{"text": "The invention relates to an expansion joint for part of a railway track disposed on a foundation which has a stock rail with a stock rail head, web and foot and a tongue movable with respect to and along said stock rail, upon which tongue at least one fastening means acts to press it onto the head of the stock rail.\nIn known expansion joints, also designated as expansion devices and permitting a movement between the structure and the rail in the vicinity of bridges, for example, a positive connection is made to permit movement of the tongue relative to the stock rail. To this end, a tongue of solid rail profile can be fixed between the stock rail and a clamming jaw disposed stationarily on the opposite side. The foot of the tongue and that of the stock rail are disposed on a common foundation at the same level. As a result, the stock rail is supported substantially on the foundation only by the stock rail half facing away from the tongue. Instabilities are compensated for by the stock rail being attached with supports and angle pieces. This entails additional maintenance work.\nA conventional rail joint is described in DE 30 16 492 A1, for example. There is a slight clearance between the clamping jaws and the facing web surface of the tongue, permitting the requisite movability of the tongue in relation to the stock rail. This too can led to tipping of the tongue or the rail."}
{"text": "The following includes information that may be useful in understanding the present invention(s). It is not an admission that any of the information provided herein is prior art, or material, to the presently described or claimed inventions, or that any publication or document that is specifically or implicitly referenced is prior art."}
{"text": "The majority of present day integrated circuits (ICs) are implemented by using a plurality of interconnected field effect transistors (FETs), also called metal oxide semiconductor field effect transistors (MOSFETs or MOS transistors). The ICs are usually formed using both P-channel and N-channel FETs and the IC is then referred to as a complementary MOS or CMOS circuit. Certain improvements in performance of FET ICs can be realized by forming the FETs in a thin layer of semiconductor material overlying an insulator layer. Such semiconductor on insulator (SOI) FETs, for example, exhibit lower junction capacitance and hence can operate at higher speeds. It is advantageous in certain applications, however, to fabricate at least some devices in the semiconductor substrate that supports the insulator layer. The devices formed in the substrate, for example, may have better thermal properties and can support higher voltages than devices formed in the thin semiconductor layer.\nAs the complexity of the integrated circuits increases, more and more MOS transistors are needed to implement the integrated circuit function. As more and more transistors are designed into the IC, it becomes important to shrink the size of individual MOS transistors so that the size of the IC remains reasonable and the IC can be reliably manufactured. Shrinking the size of an MOS transistor implies that the minimum feature size, that is, the minimum width of a line or the minimum spacing between lines, is reduced. MOS transistors have now been aggressively reduced to the point at which the gate electrode of the transistor is less than or equal to 45 nanometers (nm) in width. Methods previously used to fabricate devices in the substrate of an SOI structure, however, have not be able to achieve the same minimum feature size in substrate devices as are realized in the devices formed in the thin semiconductor layer.\nAccordingly, it is desirable to provide a method for fabricating SOI devices having minimum feature size. In addition, it is desirable to provide a self aligned method for fabricating SOI devices having minimum feature size substrate devices. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "In order to explain the background of the present invention in detail, reference will be made to FIG. 1, which shows a circuit diagram of a prior art semiconductor memory device. There are provided enhancement type MOS field-eeffect transistors 1,2,3 and 4, hereinafter referred to as MOSFETs. The drains of the P-channel MOSFET 1 and N-channel MOSFET 2 are connected to each other, and the gates thereof are connected to each other. The source of the MOSFET 1 is connected to a power supply terminal 5, and that of the MOSFET 2 is connected to ground, thus constituting a complementary MOS (hereinafter referred to as CMOS) inverter 30a. Likewise, the P-channel MOSFET 3 and the N-channel MOSFET 4 constitute a CMOS inverter 30b. With these two inverters 30a and 30b a bistable circuit, that is, a flip-flop, is formed. More particularly, the outputs of the two inverters 30a and 30b are connected to the inputs of the mating inverters 30b and 30a. In other words, the drains of the P-channel MOSFETs 1 and 3, and of the N-channel MOSFETs 2 and 4 are connected to the gates of the N-channel MOSFETs 4 and 2, and of the P-channel MOSFETs 3 and 1, respectively. In this way a one bit memory cell 30 is constructed.\nThe N-channel MOSFETs 6 and 7 which are used for transfer gates to control the writing-in and the reading-out operation, have drains (or sources) connected to the drain of the MOSFETs 1 and 2, and that of the MOSFETs 3 and 4, respectively, and have sources (or drains) connected to bit lines 8 and 9 respectively, which function as information lines for writing-in as well as reading-out. The gates of the N-channel MOSFETs 6 and 7 are connected to a word line 10 which functions as a selector line for writing-in as well as reading-out.\nThe sources and gates of the N-channel MOSFETs 11 and 12 are connected to power supply terminals 5, and their drains are connected to the bit lines 8 and 9. An information input signal line 13 is connected to the gates of the P-channel MOSFET 14 and N-channel MOSFET 15, which constitute a writing-in circuit 40. In addition, the information input signal line 13 is connected to the drain (or source) of the N-channel MOSFET 16, which is used for a gate to control the information to be written in. The source (or drain) of the MOSFET 16 is connected to the bit line 9, and its gate is connected to a writing-in control signal line 17 which is designed to control the writing operation of the memory cell 30. The drains of the MOSFETs 14 and 15 are connected to the drain (or source) of the N-channel MOSFET 18, which is used for a gate to control the data to be written in. The source (or drain) of the N-channel MOSFET 18 is connected to the bit line 8, and its gate is connected to the writing-in control signal line 17. In this way the MOSFETs 16 and 18 can transmit the output from the writing-in circuit 40 to the bit lines 8 and 9 through between the drain and source thereof.\nIn operation, the memory cells 30 and the MOSFETs 6, 7 are arrayed in matrix in plurality. A desired memory cell is directly selected by the random access method, in or from which memory cell the data is written or read out. While the memory cell stores data, the word line 10 is kept at almost zero voltage, thereby turning off the MOSFETs 6 and 7. The memory cell 30 constituted by the MOSFETs 1, 2, 3 and 4 is electrically separated from the bit lines 8 and 9. The memory cell 30 is in one of two stable states when the gates of the MOSFETs 1 and 2 are kept \"L\" (low). At this time the MOSFET 1 is in ON state with its drain being kept \"H\" (high). Accordingly, the gates of the MOSFETs 3 and 4 become \"H\", thereby turning on the MOSFET 4 with placing its drain \"L\".\nWhen the memory cell 30 is in this stable state, information can be written therein by applying voltage corresponding to the information to the bit lines 8 and 9, and applying the voltage \"H\" to the word line 10 so as to address the memory cell 30.\nNow, suppose that the logic \"1\" is to be written in the memory cell 30. The voltage \"H\" is applied to the writing-in control signal line 17, thereby turning on the MOSFETs 16 and 18, and the voltage \"H\" corresponding to the logic \"1\" is applied to the information input signal line 13. In this way the bit line 9 is kept \"H\" through the MOSFET 16. In addition, the gates of the MOSFETs 14 and 15 are kept \"H\", thereby turning off the MOSFET 14 and turning on the MOSFET 15. Thus the drains of the MOSFETs 14 and 15 become \"L\", thereby placing the bit line 8 \"L\" through the MOSFET 18.\nAt this stage, when the word line 10 is placed \"H\", the MOSFETs 6 and 7 are turned on, thereby enabling the potentials in tne bit 1ines 8 and 9 to be impressed on the memory cell 30. As a result, the MOSFET 1 is turned off whereas the MOSFET 2 is turned on, thereby reversing the states of the MOSFETs 1, 2, and 3, 4. In this way the memory cell 30 enter into the other stable state which means storing the information \"1\". Subsequently, the word line 10 and the writing-in control signal line 17 are returned to \"L\". With this, the writing operation ends.\nWhen information is to be read out from the memory cell 30, voltage of the same amplitude as that applied while writing-in operation, is impressed on the word line 10, thereby turning on the MOSFETs 6 and 7. This ensures that the electric charges stored in the bit lines 8 and 9 through the MOSFETs 11 and 12 are absorbed by the information stored in the memory cell 30, whereby a potential difference is given to between the bit lines 8 and 9 in accordance with the information stored in the memory cell 30. In this way the stored information is transmitted to the bit lines 8 and 9, and thereafter it is amplified as by a sense amplifier, and is output to the outside.\nWhen this reading operation is to be performed, it is the common practice to charge up the bit lines 8, 9 previously up to the \"H\" voltage through the MOSFETs 11, 12. This is important in preventing an erroneous writing of the information in the bit lines onto the memory cell, which is likely to occur when the MOSFETs 6, 7 are turned on in a situation where the bit lines having a large parasitic capacity have information opposite to that stored in the memory cell.\nUnder the prior art semiconductor memory device mentioned above, the bit lines are constantly charged in spite of the fact that the charging-up is required only when a reading-out operation is to be performed. As a result, the writing information and the electric charges in the bit lines come into collision when information is to be written in. This increases the consumption of electricity, and slows down the operational speeds.\nOne of prior art methods of controlling the writing and the reading operation is a technique disclosed in the article entitled \"A 4K Static 5V RAM\" by Jeffrey M. Schlageter, Nagab Jayakumar, Joseph H. Kroeger and Vahe Sarkissian, which was prepared for the 1976 International Solid-State Circuit Conference. The article teaches that by disabling the Chip Enable signal, the bit and data lines are equalized to an intermediate voltage of the power supply voltage."}
{"text": "In an inkjet recording device comprising a recording head for separately discharging ink for primary colors of cyan, magenta, and yellow and black in response to a recording signal based on image data, there are problems that (1) a secondary color of red, green, or blue is difficult to represent by adjusting primary colors and that (2) a good image cannot be obtained because of severe color mixing at the boundary between different secondary colors. To alleviate these problems, there is an inkjet recording device comprising a recording head for separately discharging ink for seven colors of cyan, magenta, yellow, black, red, green, and blue in response to a recording signal based on image data (see Patent Document 1).\nAccording to the technique disclosed in Patent Document 1, ink per se can be adjusted without complex color processing, and the secondary colors of red, green, and blue are recorded with ink of red, green, and blue, respectively. This can reduce ink injection quantity and prevent the occurrence of blur at the boundary between different colors, which would otherwise cause problems especially in secondary colors.    [Patent Document 1] JP 8-244254 A"}
{"text": "1. Field of the Invention\nThe present invention relates to a method for fabricating a FinFET transistor device, and more particularly to a method of fabricating a double gate MOSFET with a capability of inhibiting the depletion effect of the conductive gate while operating this device, which can hence elevate the device drive-on currents.\n2. Description of the Prior Art\nIn the past several years, significant progress has been made for the scaling of classical planar MOSFET (metal-oxide-semiconductor field effect transistor) structure to the gate lengths below 65 nm. Despite difficulties in fabrication, sub-20 nm physical gate length MOSFETs have recently been demonstrated. However, further scaling of the planar structure below 50 nm becomes increasingly challenging due to excess leakage, degradation in mobility, and a variety of difficulties within the device processing. Therefore, alternative MOSFET structures and new process technologies need to be explored. The effective control of leakage in nano-scale transistors will be extremely important for high-performance densely packed chips such as microprocessors.\nThe double gate MOSFET, featuring excellent short-channel behavior and relaxed requirements for aggressive scaling of gate dielectrics and junction depths, is attractive for high-performance low power applications. Fully-depleted vertical double gate devices with symmetric gate structures feature a low vertical electrical field in the channel which is favorable to carrier transport. The FinFET is a promising vertical double gate structure, which has been demonstrated in the last few years at gate lengths ranging from 100 nm to 10 nm. Unlike other reported vertical double gate structures, the FinFET can be fabricated with minimal deviation from the industry standard CMOS process.\nThe FinFET transistors were fabricated on SOT (silicon-on-insulator) wafers with a modified planar CMOS process. FIG. 1A to FIG. 1F is an illustration of the general process flow for fabricating the FinFET transistor. At first, referring to FIG. 1A, a substrate such as a silicon on insulator (SOI) structure is provided, including a silicon substrate 10, a buried oxide layer (BOX) 11 and a silicon layer 12 on the buried oxide layer 11. A cap oxide layer 13 was thermally grown on the silicon layer 12 to relieve the stress between the ensuing nitride hard mask and the silicon layer 12. A silicon nitride layer 14 was deposited on the cap oxide layer 13, to serve as a hard mask.\nAfter the hard mask deposition, a photoresistor 15 is applied to define the hard mask through use of optical lithography, electron beam lithography, X-ray lithography, or other conventional means to produce a chemical etchant mask. Then, referring to FIG. 1B, after the mask definition, an etch process is used to pattern the hard mask and the device fin structure including a silicon drain region (not shown) and a silicon source region (not shown) connected by a silicon fin or channel 12, and thereafter the photoresistor 15 is removed. Referring to FIG. 1C, a thin sacrificial oxidation process maybe used to form a sacrificial oxide 16 on the two parallelly opposing sidewalls of the silicon fin 12 to repair any damage done to the fin surface during the etch process. Oxidation may also be used to reduce the fin width, thereby allowing sub-lithography dimensions to be achieved. The threshold (Vt) implants of NMOS and PMOS can be subsequently proceeded. Referring to FIG. 1D, the hard mask of the silicon nitride layer 14 and the thin sacrificial oxide 16 are removed to retain the cap oxide layer 13 on the silicon fin 12. Referring to FIG. 1E, a gate oxide 17 is grown or deposited on the two opposing sidewalls of the silicon fin 12. Then, referring to FIG. 1F, a polysilicon gate material is deposited over the surface of the silicon fin 12, a gate mask is defined on the polysilicon gate material and then the underlying gate material is etched to form a polysilicon gate 18 with the etching stopping on the cap oxide 13 and the buried oxide layer 11. An ion implantation process is performed to implant dopants in the polysilicon gate 18 for a desired threshold voltage Vth. Referring to FIG. 2, the source/drain regions 12A and 12B are also doped to make them electrically conductive electrodes in the subsequent source/drain ion implantations. Referring to FIG. 1F again, however, the distribution of the dopants in the polysilicon gate 18 is a gaussian distribution along the depth of the polysilicon gate 18. The portions of the polysilicon gate 18 nearby the bottom corners contained between the gate oxide 17 and the buried oxide layer 11 would have been more lightly doped, and have a larger resistance. Furthermore, upon operating the FinFET transistor, a depletion of the polysilicon gate 18 easily occurs, which in turn thickens the equivalent oxide thickness (EOT) of the gate dielectric 17, resulting in the reduction of the drive current of the FinFET transistor.\nAccordingly, it is an intention to provide a method of fabricating a FinFET transistor device, which can alleviate the problem encountered in the conventional process for fabricating the FinFET transistor."}
{"text": "The present invention relates to a method for wave soldering printed circuit boards wherein a solder coating is applied only where needed.\nA usual automatic wave soldering apparatus includes a pair of endless chain conveyors driven to advance a printed circuit board at a constant speed from the entrance to exit ends of the apparatus. With the printed circuit board held by gripping fingers, the board is first carried to a fluxer where a foam or spray of flux is applied to the underside of the board. The printed circuit board is then carried over preheaters where the temperature of each board is elevated to approximately 110° C. to 130° C. so as to evaporate excess flux solvent, activate the flux and minimize thermal shock to the printed circuit board. After the printed circuit board is brought to such a preheat temperature, the board is passed over a solder reservoir to receive solder. The board is finally transported to a cool down zone where the solder is cooled to solidify.\nTypically, pin grid alley modules and dual in-line packages are loaded onto one side of a printed circuit board, with their terminals or leads projecting downwardly through apertures in the printed circuit board. To increase packaging density, a number of surface mounted devices and connectors are loaded onto the other, underside side of the printed circuit board. Problems have arisen from wave soldering such a printed circuit board. Too much heat occurs on the underside of the board, when contacted with a solder wave, and tends to damage the surface mounted devices.\nAttempts have been made to locally apply solder to preselected conductor areas on a printed circuit board, but not to those areas where surface mounted devices and connectors are mounted. In one known method, flux is locally spayed onto preselected conductor areas on a printed circuit board. The fluxed board is then preheated by a stream of warm gas to evaporate flux solvent. Thereafter, the board is positioned over a plurality of solder wave nozzles arranged within a solder reservoir. At this time, the solder wave nozzles are brought into alignment with a plurality of sets of preselected areas on the printed circuit board. A pump is arranged within the solder reservoir to force heated molten solder to flow upwardly through the solder wave nozzles so as to form solder waves. The preselected areas on the printed circuit board are contacted with the respective solder waves to make soldered joints. The height of the solder waves is then lowered until it becomes equal to the surface level of the molten solder within the solder reservoir. The position of the board is maintained until the solder solidifies. Finally, the board is delivered to a cool down zone wherein the printed circuit board is cooled.\nThis known method has proven to be effective for eliminating heat damage to sensitive electronic components such as surface mounted devices, but the method suffers from certain disadvantages. One problem is the attachment of solder oxides, better known as dross, and carbonized flux to soldered joints on a printed circuit board. Part of the molten solder tends to remain on the inner wall of the solder wave nozzles when the molten solder is immediately returned to the solder reservoir. The solder, typically composed of tin and lead, has a tendency to oxidize in the atmosphere. The resulting oxides are detrimental to the quality of the soldered joints. Also, part of the flux which has previously been applied to the soldered joints could be attached to the inner wall of the solder wave nozzles. This flux is carbonized by the heat from the molten solder. The resulting carbide could be attached to the inner wall of the solder wave nozzles. The oxides and the carbide will be separated from the inner wall of the solder wave nozzles when the molten solder within the solder reservoir is pumped up through the solder wave nozzles to process a next printed circuit board. The oxides and the carbide could be attached to soldered joints formed on the next printed circuit board. Another problem is the formation of icicles, solder bridges and other imperfections in the solder. When the conductor areas on the board are not sufficiently preheated, the solder solidifies before the conductor areas are completely wetted. This results in faulty solder connections. Also, such imperfections occurs when the solder wave is detached from the printed circuit board too fast.\nAccordingly, it is an object of the present invention to provide a method for locally applying solder to preselected conductor or solderable areas on a printed circuit board, which prevents the occurrence of icicles, bridges and other solder imperfections, and which can force solder into through holes and other hard-to-reach areas to be soldered."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device and a method of fabricating the same and, more particularly, to a semiconductor device capable of increasing an effective channel length and width, which has a plurality of grooves and a method of fabricating the same.\n2. Description of the Related Art\nAs semiconductor memory devices are highly integrated, a width of a gate electrode of a MOS transistor is reduced to about 0.1 μm and a channel length of the MOS transistor must be reduced. With reduced channel length, electrons can easily pass a gate insulating layer of the MOS transistor, and the so called “hot carrier effect” occurs because an electric field between a source and a drain of the MOS transistor is strengthen by the reduced channel length, causing flow of leakage current. Thus, the short channel effect from a reduced channel length of the MOS transistor adversely affects characteristics of the MOS transistor.\nOne way to alleviate the short channel effect is by elongating the channel length without changing a design rule of the gate electrode.\nFIGS. 1A to 1C are cross-sectional views of structures illustrating a conventional fabricating method of a semiconductor device having grooves for elongation of effective channel length.\nReferring to FIG. 1A, nitride layer patterns 12 are formed on a semiconductor substrate 10 having a device isolation layer (not shown) to expose a channel region of a MOS transistor. Next, insulating layer spacers 14 are formed on both lateral sides of the nitride layer patterns 12. Then, a thermal oxide layer 16 is formed on an exposed semiconductor substrate 10 by a thermal oxidation process.\nReferring to FIG. 1B, the thermal oxide layer 16 and the insulating spacers 14 are removed by a wet etching process to form a groove 18. Next, ions 20 are implanted on the semiconductor substrate 10 adjacent to the groove 18 for controlling a threshold voltage. Next, a gate insulating layer 22 is deposited on the semiconductor substrate 10 having the groove 18 at a predetermined thickness, and a polysilicon layer 24 is deposited on the semiconductor substrate 10 having the gate insulating layer 22 to fully fill the groove 18 between the nitride layer patterns 12. The polysilicon layer 24 is planarized to form a gate electrode 24 by a chemical mechanical polishing (CMP) process until the nitride layer pattern 12 is exposed.\nReferring to FIG. 1C, the nitride layer pattern 12 formed on both lateral sides of the gate electrode 24 is removed, and spacers 28 are formed on both lateral sides of the gate electrode 24. Conjunction regions 26 and 30 having a lightly doped drain (LDD) are formed beneath the spacers 28 and adjacent to the bottom of the spacers 28. As a result, a channel is formed under the gate electrode 24. The channel is formed to be elongated since the bottom surface of the gate electrode has a groove.\nHowever, as design rule of the gate electrode 24 is about 0.1 μm, the distance between the nitride patterns 12 should be within the design rule. During the thermal oxidation process for forming the thermal oxide layer 16, the semiconductor substrate 10 is subject to severe stress because of the nitride layer pattern 12, even if the spacers 14 are formed. Further, when the thermal oxidation process is partially performed, a bird's beak can be easily formed, thereby making difficult the forming of a delicate groove. Thus, there still continues to be a need for an improved fabricating method for elongating an effective channel length and width to improve current characteristics of a semiconductor device."}
{"text": "Production of pharmaceutically bioactive peptides and proteins in large quantities has become feasible (Biomacromolecules 2004; 5:1917-1925). The oral route is considered the most convenient way of administering drugs for patients or an animal subject. Nevertheless, the intestinal epithelium is a major barrier to the absorption of hydrophilic drugs such as peptides and proteins (J. Control. Release 1996; 39:131-138). This is because hydrophilic drugs cannot easily diffuse across the cells through the lipid-bilayer cell membranes. Attentions have been given to improving paracellular transport of hydrophilic drugs (J. Control. Release 1998; 51:35-46). However, the transport of hydrophilic molecules via the paracellular pathway is, however, severely restricted by the presence of tight junctions that are located at the luminal aspect of adjacent epithelial cells (Annu. Rev. Nutr. 1995; 15:35-55). These tight junctions form a barrier that limits the paracellular diffusion of hydrophilic molecules. The structure and function of tight junctions is described, inter alia, in Ann. Rev. Physiol. 1998; 60:121-160 and in Ballard T S et al., Annu. Rev. Nutr. 1995; 15:35-55. Tight junctions do not form a rigid barrier but play an important role in the diffusion through the intestinal epithelium from lumen to bloodstream and vice versa.\nMovement of solutes between cells, through the tight junctions that bind cells together into a layer such as the epithelial cells of the gastrointestinal tract, is termed paracellular transport. Paracellular transport is passive. Paracellular transport is dependent on electrochemical gradients generated by transcellular transport and solvent drag through tight junctions. Tight junctions form an intercellular barrier that separates the apical and basolateral fluid compartments of a cell layer. Movement of a solute through a tight junction from apical to basolateral compartments depends on the permeability of the tight junction for that solute.\nPolymeric nanoparticles have been widely investigated as carriers for drug delivery (Biomaterials 2002; 23:3193-3201). Much attention has been given to the nanoparticles made of synthetic biodegradable polymers such as poly-ε-caprolactone and polylactide due to their biocompatibility (J. Drug Delivery 2000; 7:215-232; Eur. J. Pharm. Biopharm. 1995; 41:19-25). However, these nanoparticles are not ideal carriers for hydrophilic drugs because of their hydrophobic property. Some aspects of the invention relate to a novel nanoparticle system, composed of hydrophilic chitosan and poly(glutamic acid) hydrogels; the nanoparticles are prepared by a simple ionic-gelation method. This technique is promising as the nanoparticles are prepared under mild conditions without using harmful solvents. It is known that organic solvents may cause degradation of peptide or protein drugs that are unstable and sensitive to their environments (J. Control. Release 2001; 73:279-291).\nFollowing the oral drug delivery route, protein drugs are readily degraded by the low pH of gastric medium in the stomach. The absorption of protein drugs following oral administration is challenging due to their high molecular weight, hydrophilicity, and susceptibility to enzymatic inactivation. Protein drugs at the intestinal epithelium cannot partition into the hydrophobic membrane, leaving only the epithelial barrier via the paracellular pathway. However, the tight junction forms a barrier that limits the paracellular diffusion of hydrophilic molecules.\nChitosan (CS), a cationic polysaccharide, is generally derived from chitin by alkaline deacetylation (J. Control. Release 2004; 96:285-300). It was reported from literature that CS is non-toxic and soft-tissue compatible (Biomacromolecules 2004; 5:1917-1925; Biomacromolecules 2004; 5:828-833). Additionally, it is known that CS has a special property of adhering to the mucosal surface and transiently opening the tight junctions between epithelial cells (Pharm. Res. 1994; 11:1358-1361). Most commercially available CSs have a quite large molecular weight (MW) and need to be dissolved in an acetic acid solution at a pH value of approximately 4.0 or lower, which is somewhat impractical. However, there are potential applications of CS in which a low MW would be essential. Given a low MW, the polycationic characteristic of CS can be used together with a good solubility at a pH value close to physiological ranges (Eur. J. Pharm. Biopharm. 2004; 57:101-105). Loading of peptide or protein drugs at physiological pH ranges would preserve their bioactivity. On this basis, a low-MW CS, obtained by depolymerizing a commercially available CS using cellulase, is disclosed herein to prepare nanoparticles of the present invention.\nThanou et al. reported chitosan and its derivatives as intestinal absorption enhancers (Adv Drug Deliv Rev 2001; 50:S91-S101). Chitosan, when protonated at an acidic pH, is able to increase the paracellular permeability of peptide drugs across mucosal epithelia. Co-administration of chitosan or trimethyl chitosan chloride with peptide drugs were found to substantially increase the bioavailability of the peptide in animals compared with administrations without the chitosan component.\nThe γ-PGA, an anionic peptide, is a natural compound produced as capsular substance or as slime by members of the genus Bacillus (Crit. Rev. Biotechnol. 2001; 21:219-232). γ-PGA is unique in that it is composed of naturally occurring L-glutamic acid linked together through amide bonds. It is reported from literature that this naturally occurring γ-PGA is a water-soluble, biodegradable, and non-toxic polymer. A polyamino carboxylic acid (complexone), such as diethylene triamine pentaacetic acid, has showed enzyme resistant property. It is clinical beneficial to incorporate a PGA-complexone conjugate as a negative substrate and chitosan as a positive substrate in a drug delivery nanoparticle formulation for better absorption performance with reduced enzymatic effect."}
{"text": "FIG. 1 shows a prior art six transistor static read/write memory cell 710 such as is typically used in high-density static random access memories (SRAMs). A static memory cell is characterized by operation in one of two mutually-exclusive and self-maintaining operating states. Each operating state defines one of the two possible binary bit values, zero or one. A static memory cell typically has an output which reflects the operating state of the memory cell. Such an output produces a \"high\" voltage to indicate a \"set\" operating state. The memory cell output produces a \"low\" voltage to indicate a \"reset\" operating state. A low or reset output voltage usually represents a binary value of zero, while a high or set output voltage represents a binary value of one.\nStatic memory cell 710 generally comprises first and second inverters 712 and 714 which are cross-coupled to form a bistable flip-flop. Inverters 712 and 714 are formed by n-channel driver transistors 716 and 717, and p-channel load transistors 718 and 719. Driver transistors 716 and 717 are typically n-channel metal oxide silicon field effect transistors (MOSFETs) formed in an underlying silicon semiconductor substrate. P-channel transistors 718 and 719 are typically thin film transistors formed above the driver transistors.\nThe source regions of driver transistors 716 and 717 are tied to a low reference or circuit supply voltage, labelled V.sub.ss and typically referred to as \"ground.\" Load transistors 718 and 719 are connected in series between a high reference or circuit supply voltage, labelled V.sub.cc, and the drains of the corresponding driver transistors 716 and 717. The gates of load transistors 718 and 719 are connected to the gates of the corresponding driver transistors 716 and 717.\nInverter 712 has an inverter output 720 formed by the drain of driver transistor 716. Similarly, inverter 714 has an inverter output 722 formed by the drain of driver transistor 717. Inverter 712 has an inverter input 724 formed by the gate of driver transistor 716. Inverter 714 has an inverter input 726 formed by the gate of driver transistor 717.\nThe inputs and outputs of inverters 712 and 714 are cross-coupled to form a flip-flop having a pair of complementary two-state outputs. Specifically, inverter output 720 is cross-coupled to inverter input 726, and inverter output 722 is cross-coupled to inverter input 724. In this configuration, inverter outputs 720 and 722 form the complementary two-state outputs of the flip-flop.\nA memory flip-flop such as that described typically forms one memory element of an integrated array of static memory elements. A plurality of access transistors, such as access transistors 730 and 732, are used to selectively address and access individual memory elements within the array. Access transistor 730 has one active terminal connected to cross-coupled inverter output 720. Access transistor 732 has one active terminal connected to cross-coupled inverter output 722. A plurality of complementary column line pairs, such as the single pair of complementary column lines 734 and 736 shown, are connected to the remaining active terminals of access transistors 730 and 732, respectively. A row line 738 is connected to the gates of access transistors 730 and 732.\nReading static memory cell 710 involves activating row line 738 to connect inverter outputs 720 and 722 to column lines 734 and 736. Writing to static memory cell 710 involves first placing selected complementary logic voltages on column lines 734 and 736, and then activating row line 738 to connect those logic voltages to inverter outputs 720 and 722. This forces the outputs to the selected logic voltages, which will be maintained as long as power is supplied to the memory cell, or until the memory cell is reprogrammed.\nFIG. 2 shows an alternative four transistor, dual wordline, prior art static read/write memory cell 750 such as is typically used in high density static random access memories. Static memory cell 750 comprises n-channel pull-down (driver) transistors 780 and 782 having drains respectively connected to pull-up load elements or resistors 784 and 786. Transistors 780 and 782 are typically metal oxide silicon field effect transistors (MOSFETs) formed in an underlying silicon semiconductor substrate.\nThe source regions of transistors 780 and 782 are tied to a low reference or circuit supply voltage, labelled V.sub.ss and typically referred to as \"ground.\" Resistors 784 and 786 are respectively connected in series between a high reference or circuit supply voltage, labelled V.sub.cc, and the drains of the corresponding transistors 780 and 782. The drain of transistor 782 is connected to the gate of transistor 780 by line 776, and the drain of transistor 780 is connected to the gate of transistor 782 by line 774 to form a flip-flop having a pair of complementary two-state outputs.\nA memory flip-flop, such as that described above in connection with FIG. 2, typically forms one memory element of an integrated array of static memory elements. A plurality of access transistors, such as access transistors 790 and 792, are used to selectively address and access individual memory elements within the array. Access transistor 790 has one active terminal connected to the drain of transistor 780. Access transistor 792 has one active terminal connected to the drain of transistor 782. A plurality of complementary column line pairs, such as the single pair of complementary column lines 752 and 754 shown, are connected to the remaining active terminals of access transistors 790 and 792, respectively. A row line 756 is connected to the gates of access transistors 790 and 792.\nReading static memory cell 750 involves activating row line 756 to connect outputs 768 and 772 to column lines 752 and 754. Writing to static memory cell 750 involves first placing selected complementary logic voltages on column lines 752 and 754, and then activating row line 756 to connect those logic voltages to outputs 768 and 772. This forces the outputs to the selected logic voltages, which will be maintained as long as power is supplied to the memory cell, or until the memory cell is reprogrammed.\nA static memory cell is said to be bistable because it has two stable or self-maintaining operating states, corresponding to two different output voltages. Without external stimuli, a static memory cell will operate continuously in a single one of its two operating states. It has internal feedback to maintain a stable output voltage, corresponding to the operating state of the memory cell, as long as the memory cell receives power.\nThe two possible output voltages produced by a static memory cell correspond generally to upper and lower circuit supply voltages. Intermediate output voltages, between the upper and lower circuit supply voltages, generally do not occur except for during brief periods of memory cell power-up and during transitions from one operating state to the other operating state.\nThe operation of a static memory cell is in contrast to other types of memory cells such as dynamic cells which do not have stable operating states. A dynamic memory cell can be programmed to store a voltage which represents one of two binary values, but requires periodic reprogramming or \"refreshing\" to maintain this voltage for more than very short time periods.\nA dynamic memory cell has no internal feedback to maintain a stable output voltage. Without refreshing, the output of a dynamic memory cell will drift toward intermediate or indeterminate voltages, resulting in loss of data. Dynamic memory cells are used in spite of this limitation because of the significantly greater packaging densities which can be attained. For instance, a dynamic memory cell can be fabricated with a single MOSFET transistor, rather than the four or more transistors typically required in a static memory cell. Because of the significantly different architectural arrangements and functional requirements of static and dynamic memory cells and circuits, static memory design has developed along generally different paths than has the design of dynamic memories.\nFIG. 3 illustrates a typical top view of a prior art layout of portions of the FIG. 2 SRAM schematic pertinent to this disclosure. Such an SRAM cell employs two Vcc lines which are labeled respectively as Vcc(A) and Vcc(B). Lines 790a and 792a constitute the gate or wordlines of access devices 790 and 792, respectively. The two shaded areas 784 and 786 comprise the described pull-up resistors which are substantially horizontally formed as shown within the SRAM cell. Such horizontal positioning consumes considerable area within an individual SRAM cell, thus countering a desired goal of maximizing circuit density. Other examples of such similarly situated resistors can be found in U.S. Pat. No. 4,178,674 to Liu et al. and U.S. Pat. No. 4,828,629 to Akeda et al. Vertically oriented pull-up resistor constructions in SRAM cells have been proposed, such as is disclosed in our U.S. Pat. No. 5,177,030.\nThis invention concerns improved SRAM construction employing vertically elongated pull-up resistors in SRAMs."}
{"text": "A DRAM is a commonly used semiconductor device comprising a capacitor and a transistor. A continuous challenge in the semiconductor industry is to decrease the vertical and/or horizontal size of semiconductor devices, such as DRAMs and capacitors. A limitation on the minimal horizontal footprint of capacitor constructions is impacted by the resolution of a photolithographic etch during fabrication of the capacitor constructions. Although this resolution is generally improving, at any given time there is a minimum photolithographic feature dimension of which a fabrication process is capable. It would be desirable to form capacitors at least some portions of which have a cross-sectional minimum dimension of less than the minimum capable photolithographic feature dimension of a given fabrication process.\nA problem in the semiconductor industry is mask misalignment. Mask misalignment during device fabrication can lead to inoperative devices. Accordingly, it is desirable to design device-fabrication processes which can compensate for mask misalignment."}
{"text": "In the broadband communication system such as a TV tuner, a highly linear variable-gain low noise amplifier is arranged in upstream of a mixer. Generally, the highly linear variable-gain low noise amplifier is implemented according to a current steering topology.\nFIG. 1A is a schematic circuit diagram illustrating a conventional highly linear variable-gain amplifier. This highly linear variable-gain amplifier is disclosed in IEEE J. Solid-State Circuits, vol. 26, pp. 1673-1680, November 1991. As shown in FIG. 1A, a first transistor Q1 and a second transistor Q2 are connected with each other to define a differential pair. The bases of the first transistor Q1 and the second transistor Q2 serve as the differential signal input terminals of the amplifier to receive an input signal vi. The first terminals of two emitter resistors Re are respectively connected to the emitters of the first transistor Q1 and the second transistor Q2. The second terminals of the two emitter resistors Re are collectively connected to a node “a”. A current source (Is) is interconnected between the node “a” and a ground terminal Gnd.\nThe bases of a third transistor Q3 and a fourth transistor Q4 serve as the gain control terminals of the amplifier for receiving a current steering control signal Vctrl. The collector of the third transistor Q3 is connected to a voltage source Vcc. The emitter of the third transistor Q3 is connected to the collector of the first transistor Q1. A first collector resistor Rc1 is interconnected between the voltage source Vcc and the collector of the fourth transistor Q4. The emitter of the fourth transistor Q4 is connected to the collector of the first transistor Q1. The base of a fifth transistor Q5 is connected to the base of the fourth transistor Q4. The base of a sixth transistor Q6 is connected to the base of the third transistor Q3. The collector of the sixth transistor Q6 is connected to the voltage source Vcc. The emitter of the sixth transistor Q6 is connected to the collector of the second transistor Q2. A second collector resistor Rc2 is interconnected between the collector of the fifth transistor Q5 and the voltage source Vcc. The emitter of the fifth transistor Q5 is connected to the collector of the second transistor Q2. The collectors of the fourth transistor Q4 and the fifth transistor Q5 serve as differential signal output terminals of the amplifier for generating an output signal vo.\nThe current source (Is) may provide DC bias voltages to all transistors of the amplifier. The two emitter resistors Re may offer good linearity of the amplifier. In addition, the resistance of the first collector resistor Rc1 is identical to that of the second collector resistor Rc2.\nIn response to a change of the current steering control signal Vctrl, the bias currents flowing through the third transistor Q3, the fourth transistor Q4, the fifth transistor Q5 and the sixth transistor Q6 are varied, and thus the gain value of the amplifier are adjustable. Moreover, the above amplifier may acquire a high gain control range.\nGenerally, the noise figure (NF) of the highly linear variable-gain amplifier is varied with the gain value. FIGS. 1B and 1C are schematic diagrams illustrating the relationship between the gain and the noise figure (NF) of the conventional highly linear variable-gain amplifier. As can be seen from FIGS. 1B and 1C, as the gain value of the amplifier is increased, the noise figure is decreased. Whereas, as the gain value of the amplifier is decreased, the noise figure is increased. That, when the gain value of the amplifier is adjusted according to the current steering control signal Vctrl, the noise figure is increased at nearly the same rate as the gain value is decreased.\nIn a case that the magnitude of the input signal vi is very low, the gain value of the amplifier is usually adjusted to the maximum value, and thus the noise figure is not too large. In a case that the magnitude of the input signal vi is relatively larger, the gain value needs to be reduced. In this situation, the noise figure of the amplifier is increased, and the signal is also amplified. In other words, the magnitude of the output signal vo allows for providing a sufficient signal-to-noise ratio (SNR). However, in the broadband communication application, the interference and noise are sometimes greater than the useful signal. For preventing the electronic components of the amplifier from entering the saturation region, the gain value of the amplifier needs to be decreased. If the increase of the noise figure is too obvious, however, the magnitude of the output signal vo fails to provide a sufficient signal-to-noise ratio (SNR), and thus the baseband circuit fails to effectively restore the signal. That is, when the amplifier has a low gain, low noise figure (NF) is very critical.\nFIG. 2 is a schematic circuit diagram illustrating another conventional highly linear variable-gain low noise amplifier. The highly linear variable-gain low noise amplifier is disclosed in for example U.S. Pat. No. 6,100,761. As shown in FIG. 2, a first transistor 1Q1 and a second transistor 1Q2 are connected with each other to define a differential pair. The base of the first transistor 1Q1 is connected with a base voltage Vb through a first base resistor 1Rb1. The base of the second transistor 1Q2 is connected with the base voltage Vb through a second base resistor 1Rb2. The bases of the first transistor Q1 and the second transistor Q2 serve as the differential signal input terminals (IN+ and IN−) of the amplifier.\nThe first terminals of two variable emitter resistors (1Re) 40 are respectively connected to the emitters of the first transistor 1Q1 and the second transistor 1Q2. The second terminals of two variable collector resistors (1Rc) 30 are respectively connected to a collector voltage Vc. Moreover, the collectors of the first transistor Q1 and the second transistor Q2 serve as the differential signal output terminals (− OUT +) of the amplifier.\nIn the amplifier of FIG. 2, the gain value of the amplifier is adjusted by changing the resistances of the variable emitter resistors (1Re) 40 and the variable collector resistors (1Rc) 30. The changes of the variable emitter resistors (1Re) 40, however, may deteriorate the linearity of the amplifier."}
{"text": "1. Field\nEmbodiments of the disclosure relate generally to the field of aircraft flight instrumentation and more particularly to a system for notifying pilots of an off ground flight condition resulting from a landing bounce which incorporates landing gear extension/compression measurement and pilot indicators.\n2. Background\nAircraft operating in high winds or other adverse environmental conditions during landing or landing with improperly serviced landing gear shock struts may encounter the runway surface with sufficient force to induce a bounce (rebound off the runway) causing the aircraft to again become airborne. Premature rapid de-rotation (pitching the nose down) without the main landing gear firmly planted on the runway has been cited as one of the contributing elements for aircraft undesired landing incidents. Due to the size and cockpit positioning in large commercial aircraft, if an airplane bounces off the runway after touchdown, the flight crew may not recognize that the airplane has bounced due to the difference between the cockpit movement and the airplane main landing gear position relative to the ground; i.e. the cockpit may be moving up or down while the main gear are coming off the ground. If the airplane derotation is performed while the main gear are off the ground, it could lead to a nose landing gear only landing that could cause significant structural damage to the airplane.\nIt is therefore desirable to provide an Airplane Off Ground Advisory System (AOGAS) which will provide an indication to the flight crew positively identifying the aircraft has bounced after initial contact with the runway so that the crew can take appropriate action to make a safe landing or perform a go-around operation."}
{"text": "MIMO (multi-input multi-output) technology corresponds to a technology for increasing data transmission and reception efficiency using a plurality of transmission antennas and a plurality of reception antennas instead of using a single transmission antenna and a single reception antenna. If a single antenna is used, a receiving end receives data through a single antenna path. On the contrary, if multiple antennas are used, the receiving end receives data through several paths, thereby enhancing transmission speed and transmission capacity and increasing coverage.\nA single-cell MIMO operation can be divided into a single user-MIMO (SU-MIMO) scheme that a single user equipment (UE) receives a downlink signal in a single cell and a multi user-MIMO (MU-MIMO) scheme that two or more UEs receive a downlink signal in a single cell.\nChannel estimation corresponds to a procedure of restoring a received signal by compensating a distortion of the signal distorted by fading. In this case, the fading corresponds to a phenomenon of rapidly changing strength of a signal due to multi-path time delay in wireless communication system environment. In order to perform the channel estimation, it is necessary to have a reference signal known to both a transmitter and a receiver. The reference signal can be simply referred to as an RS (reference signal) or a pilot depending on a standard applied thereto.\nA downlink reference signal corresponds to a pilot signal for coherently demodulating PDSCH (physical downlink shared channel). PCFICH (physical control format indicator channel), PHICH (physical hybrid indicator channel). PDCCH (physical downlink control channel) and the like. A downlink reference signal can be classified into a common reference signal (CRS) shared by all UEs within a cell and a dedicated reference signal (DRS) used for a specific UE only. Compared to a legacy communication system supporting 4 transmission antennas (e.g., a system according to LTE release 8 or 9 standard), a system including an extended antenna configuration (e.g., a system according to LTE-A standard supporting 8 transmission antennas) is considering DRS-based data demodulation to efficiently manage a reference signal and support an enhanced transmission scheme. In particular, in order to support data transmission through an extended antenna, it may be able to define a DRS for two or more layers. Since a DRS and data are precoded by a same precoder, it is able to easily estimate channel information, which is used for a receiving end to demodulate data, without separate precoding information.\nAlthough a downlink receiving end is able to obtain precoded channel information on an extended antenna configuration through a DRS, it is required for the downlink receiving end to have a separate reference signal except the DRS to obtain channel information which is not precoded. Hence, it is able to define a reference signal for obtaining channel state information (CSI), i.e., a CSI-RS, at a receiving end in a system according to LTE-A standard."}
{"text": "1. Field of the Invention\nThis invention relates to frequency and award redemption program. More particularly, the present invention relates to an on-line, interactive frequency and award redemption program which is fully integrated.\n2. Description of Related Art\nFrequency programs have been developed by the travel industry to promote customer loyalty. An example of such a program is a “frequent flyer” program. According to such a program, when a traveler books a flight, a certain amount of “mileage points” are calculated by a formula using the distance of the destination as a parameter. However, the mileage points are not awarded until the traveler actually takes the flight.\nWhen a traveler has accumulated sufficient number of mileage points, he may redeem these points for an award chosen from a specific list of awards specified by the program. Thus, for example, the traveler may redeem the points for a free flight ticket or a free rental car. In order to redeem the accumulated points, the traveler generally needs to request a certificate, and use the issued certificate as payment for the free travel.\nWhile the above program may induce customer loyalty, it has the disadvantage that the selection of prizes can be made only from the limited list of awards provided by the company. For example, a traveler may redeem the certificate for flights between only those destinations to which the carrier has a regular service. Another disadvantage is that the customer generally needs to plan ahead in sufficient time to order and receive the award certificate.\nAccording to another type of frequency and award program, a credit instrument is provided and credit points are accumulated instead of the mileage points. In such programs, bonus points are awarded by using a formula in which a price paid for merchandise is a parameter. Thus, upon each purchase a certain number of bonus points are awarded, which translate to dollar credit amount. According to these programs, the customer receives a credit instrument which may be acceptable by many enrolled retailers, so that the selection of prizes available is enhanced. An example of such a program is disclosed in E.P.A. 308,224. However, while such programs may enhance the selection of prizes, there is still the problem of obtaining the credit instrument for redeeming the awarded points. In addition, the enrollee must allow for processing time before the bonus points are recorded and made available as redeemable credit. Thus, the immediacy effect of the reward is lacking in these conventional incentive programs."}
{"text": "U.S. Pat. No. 5,306,305 (=WO 95/19152) discloses an in vitro method for producing an implant device by coating a gel containing osteoblast cells onto a porous metal surface and then incubating the gel in a growth medium. A repeatedly renewed minimal essential medium (MEM) is used for about 3 weeks for cell multiplication, followed by a medium containing .beta.-glycerophosphate and ascorbic acid for another 1-2 weeks. The cells may originate from the patient's own bone fragments. The gel (e.g. 0.5% gelatin) is used to hold the cells to the substrate surface.\nDE-A-3810803 discloses a method of producing living bone substitute materials by in vitro culturing autologous bone cells from human bone fragments in a repeatedly renewed culture medium, followed by deposition of the cultured cells in a porous calcium phosphate matrix and additional culturing. The composite material can be reimplanted.\nWO 94/04657 discloses a bioactive porous glass which is pretreated in such a way that it cannot raise the pH of a tissue medium contacted with the glass. It also reports the seeding of the pretreated porous glass with osteoblasts.\nThese prior art methods of in vitro production of bone tissue for implanting purposes have not yet been put into practice, probably because fixation of the resulting implant in the body and thus functioning of the implant are insufficient due to limitations in the applying techniques. No biological effect of using a particular culture method was described in the prior art. Furthermore, these prior methods necessitate the introduction of a bone defect (a lesion) in the patient in order to obtain the required bone cells."}
{"text": "The present invention relates generally to an apparatus for the distribution and delivery of video in a digital system and, in particular, to a method and apparatus that rapidly changes the channel in such a system.\nIn a typical analog cable system all of the channels or services ordered by the subscriber are delivered to each subscriber\"\"s home. In order to ensure that each subscriber receives only the channels for which he has paid, the cable television providers encrypt or xe2x80x9cscramblexe2x80x9d the premium channels (HBO, CINEMAX, DISNEY, etc.). The cable television providers also may scramble many of the xe2x80x9cbasicxe2x80x9d channels (local stations, ESPN, MTV, VH1, TNT, DISCOVERY, etc.). Therefore, even though virtually all television sets sold today are cable-ready, most subscribers still need a set-top unit (sometimes referred to as a cable box in the cable television environment) to descramble the signals. The set-top units are located proximate a television and are also used to change the channel that is viewed on the television.\nSubscribers, especially residential subscribers, are demanding that large amounts of information and more choices of services be brought into their homes. Switched video for viewing on a subscriber\"\"s television and high-speed Internet access are two services highly desired by subscribers. In order to meet the demand, and to accommodate the recently approved high-definition television standards, the new services will likely have to be capable of handling digital signals.\nIn addition to the emerging technologies, entry of the Regional Bell Operating Companies (RBOCs) have made digital delivery systems economically feasible. A typical digital video delivery system includes a means for receiving the video signals from various broadcast sources, a means for delivering the signals to a plurality of subscribers, and a means of transmitting the signals between the receiving means and delivering means. The means for receiving the video signals may include a broadband digital terminal (BDT) located in a central office. The delivering means may be a broadband network unit (BNU) located preferably on a telephone pole or other convenient location proximate a number of subscribers. Cable or optical fiber connects the BDT to the BNU. A second cable (or a twisted wire pair) connects the set-top units (and, if required, the various other units in the subscriber\"\"s home) with the BNU.\nIn a typical digital video system, all of the video services offered by the video service company are again delivered to the set-top unit. When selecting a new channel, the set-top unit performs the actual switching and also descrambles the digital signals.\nEven though the premium channels are scrambled, delivering all of the video signals into a subscriber\"\"s home makes the video services susceptible to theft. Accordingly, more complexxe2x80x94and expensivexe2x80x94steps must be taken to further secure the transmission and delivery of the video services.\nA solution to the theft problem is to perform the channel switching xe2x80x9cupstreamxe2x80x9d from the subscriber at a facility controlled by the video provider (in a digital system at the BDT, for, example), and only delivering one channel at a time to the subscriber\"\"s set-top unit. Another advantage of moving the switching upstream is that the bandwidth requirements of the overall video delivery system are greatly reduced. However, a drawback of this system is that the subscriber experiences a time delay between the period of time it takes for the subscriber to select a channel and for the newly selected channel to be viewed on the television.\nThe reason for this time delay is that the subscriber\"\"s request must first travel upstream to the BDT; next, the BDT must acknowledge the request, then synchronize and xe2x80x9clock onxe2x80x9d to the desired video service; finally, the BDT must transmit the desired video service back downstream to the subscriber. The overall delay between each channel change can take over a second.\nA significant portion of the delay is caused by the time it takes for the video signal to synchronize. This portion accounts for about half of the overall delay (i.e., about a half second). Many subscribers find the delay in such video delivery systems annoying since they are accustomed to seeing the broadcast signal immediately after selecting a new channel. Most subscribers find that a one second delay is unacceptable, which would make such a system competitively unattractive.\nThe present invention relates to a method and apparatus for rapidly changing the channel in a digital video delivery system. The rapid channel changer will be preferably located in the central office with the broadband digital terminal (BDT) and indexes the xe2x80x9cstartxe2x80x9d or synchronization frame of each video channel received at the BDT.\nEach digital video signal includes a synchronization frame. The subject channel changer captures the multiple compressed video signals and stores each signal in a cache buffer. A, processor is used to index or xe2x80x9cpoint toxe2x80x9d the respective synchronization frames for each buffered signal.\nWhen a subscriber requests a specific channel or video service, the processor can immediately access the requested video signal at a synchronization frame and direct the video stream to the subscriber since the processor already has the position of the synchronization frame of each video signal. Accordingly, the period of time that the subscriber previously had to wait for the synchronization frame is eliminated.\nThese and other features and objects of the invention will be more fully understood from the following detailed description of the preferred embodiments which should be read in light of the accompanying drawings."}
{"text": "In wafer inspection systems which utilize two dimensional imaging, the inspection speed is determined, among other things, from parameters including field of view size, and time between imaging sequential images. Generally speaking, a larger field of view, or a shorter time between sequential images will increase the inspection speed.\nDecreasing the time between imaging may be complicated and expensive. For instance, decreasing the time between images can require very fast detectors (much faster above normal 30 Hz detectors), fast illumination (for example, repetitive laser with hundreds of pulses per second), and a fast stage or other suitable components for generating relative motion between the wafer and imaging components to change which portion(s) of the wafer are in view for imaging.\nA more preferable approach in some circumstances is to enlarge the field of view. However, when fine resolution is required (pixel size in the wafer plane is below 0.5 microns), the detector must contain a numerous pixels. For example, using 0.2 micron pixel, and a conventional commercial detector with 2K×2K pixels, the field of view is only 0.4 mm×0.4 mm. An enlarged field of view may also require a faster stage or other suitable components for providing relative motion between the imaging components and the wafer.\nThe image view can be increased by using multiple two dimensional detectors to obtain an image, with the image divided amongst the detectors. Some currently-existing systems split an image before the focal plane of the other optics used to obtain the image using, for instance, beam splitters and/or mirrors. See, for instance, U.S. patent application Ser. No. 10/345,097, filed Jan. 15, 2003, and published as U.S. Patent Application Publication No., 20040146295 which are each incorporated by reference in their entireties herein. However, splitting an image by a mirror or other element(s) before the focal plane may be problematic in some instances. The problems may include, for example, reductions in intensity and/or non-uniform intensity.\nFIG. 9 illustrates an example wherein the intensity in some parts of a split image is reduced when some rays are reflected back from the mirror and do not actually reach the focal plane, since the actual splitting of the image occurs prior to the focal plane. As shown in FIG. 9, three rays (R1, R2 and R3) from the imaging optics 18 of an inspection system reach point A in the focal plane FP18 of the imaging optics if no splitting mirror is used (i.e., if the mirror shown in FIG. 9 is disregarded, all three rays reach point A). However, when the splitting mirror comprising reflective planes 902 and 904 is used, only two rays (R2 and R3) reach the detectors 908-1 and 908-2 in the split focal plane. The top ray (R1) is reflected back from the mirror.\nFIG. 9 also illustrates an example of non uniform intensity that may result from splitting. The intensity reduction is position dependent—a given portion of the image that is closer to the splitting point will have a reduced intensity relative to a portion of the image far from the splitting point. In FIG. 9, point B′ gets only about half of the rays (i.e. rays generally emanating from the bottom half part of the imaging optics), while point A′, for example, gets more (about two thirds: from ray R2 to R3).\nAn example hypothetical intensity distribution in detector 908-1 and 908-2 imaging a uniform input image (I and II) is shown in FIG. 10. The image is darker at points closer to the splitting point, with denser cross-hatching representing progressively darker portions of the image (becoming darker from left to right in 908-2 and right to left in 908-1).\nThe angular distribution of the image is not preserved when an image is split in this manner. For a wafer inspection system, the angular distribution of the scattered or reflected light from the wafer contains information regarding the wafer characteristics. Using splitting mirrors before the focal plane changes the angular distribution since it blocks a range of ray angles and thus may result in reduced inspection accuracy.\nWhen splitting by beam splitters, some of the rays (usually 50%) are reflected from the beam splitter while the rest of the rays are transmitted. This way does not break the uniformity or the angular distribution, but the intensity is reduced by 50%. When using more than one splitter to split an image into more than two portions, the intensity can be reduced even more."}
{"text": "1. Field of the Invention\nThe present invention relates to a display panel drive circuit and display panel and, more particularly, to a display panel drive circuit and display panel in which thin film transistors for the display panel drive circuit can be prevented from deteriorating.\n2. Description of Related Art\nIn recent years, there have been proposals for LCD (liquid crystal display) panels utilizing low-temperature polysilicon TFTs (thin film transistors). Such display panels can be formed, on one common substrate, not only together with pixel transistors but also with peripheral drive circuits, such as scanning shift registers and sampling circuits. Accordingly, display can be by mere external connection with reduced number of signal lines, reducing the number of parts and improving reliability. Large display panels of an approximately 20-40 type are under considerations.\nThere are recently found cases where a large color liquid crystal panels is equipped on a camera-integrated video tape recorder (VTR) in order for use as monitors or finders. Of these camera-integrated VTRs, there are structures that a display panel is arranged to rotate about a horizontal axis to shift its position. In such a case, horizontal and vertical scanning directions has to be changed depending upon the panel direction so that display is properly viewed when the panel is rotated. Due to this, the scanning shift register includes a scanning direction control circuit using, for example, an analog switch circuit.\nA display panel having a drive circuit for controlling the scanning direction, as mentioned above, was formed on one substrate, for conducting test. It was confirmed that deterioration is encountered in the TFTs of a signal input circuit to which scanning start pulses are externally applied, causing a problem that initial failure occurs resulting in impossibility of scanning.\nThe cause of such deterioration in the signal input circuit TFTs externally applied by scanning start pulses is to be presumed as follows. That is, the start pulse drive circuit is high in deriveability, and circuit board mounting is separated from the display panel with connections to the display panel through cables, flexible circuit boards or the like. During driving or upon switching the scanning direction, a high voltage occurs due to the effect of interconnection inductance, etc., resulting in deterioration or breakage of transistors. It is also to be presumed as one of reasons for the deterioration that the analog switch circuit, to which an external start pulse is first inputted, is not configured as a gate input circuit.\nAlso, where the panel is made larger, a problem of time delay occurs particularly for a pixel section. In such a case, there is a necessity of forming the interconnection (gate) with using a low-resistance material such as aluminum. In the above-stated display panel, however, a pixel section and its peripheral circuits are formed by a common process so that the interconnections for the peripheral circuit are formed also by the low-resistance material. Due to this, there has been a problem that the peripheral circuit elements are liable to undergo dielectric breakdown.\nFurther, where using a high insulation substrate such as a glass substrate, there occurs concentration of electric fields through the interconnections during a plasma process for the TFT manufacture, resulting in a problem that so-called plasma antenna effects occur, i.e., the elements connected to these interconnections undergo damage. This phenomenon is liable to occur, particularly, at end portions of an interconnection pattern, at discontinuous portions or at large electrode areas. This condition is met by a start pulse input terminal pattern.\nIt is an object of the present invention to solve the above-stated problems as encountered in the prior art, and to provide a display panel drive circuit and display panel which is simple in structure but is free from occurrence of initial failure leading to impossibility of scanning.\nA display panel drive circuit according to the present invention is characterized in that: thin film transistors constituting a signal input circuit connected to a circuit outside the display panel are formed in a structure having a dielectric breakdown strength higher than those of thin film transistors constituting other circuits.\nIn the present invention, only a circuit to which signals are externally applied or thin film transistors of the same circuit is formed by a structure to withstand high voltage, whereby they operate in a manner preventing against deterioration and hence occurrence of initial failure."}
{"text": "(Not Applicable)\n(Not Applicable)\nThe present invention generally relates to an illumination head for printed circuit board verification and more particularly to an illumination head that may be disposed a prescribed distance above the printed circuit board in order to facilitate rework thereon.\nIn order to inspect printed circuit boards, it is desirable to illuminate the same with an illumination source. Typical illumination sources include rings of light and/or spotlights which direct light onto a top surface of the printed circuit board. In this regard, once the top surface of the printed circuit board is illuminated, the printed circuit board may be inspected through the use of a video camera and/or automated inspection technique.\nIt is desirable to illuminate the top surface of the printed circuit board such that shadows and/or glare which can impair inspection are eliminated. Therefore, in the prior art, the illumination head for an inspection station is disposed in close proximity to the top surface of the printed circuit board.\nIt is undesirable to have the illumination head disposed in close proximity to the top surface of the printed circuit board because reworking of a defective circuit board is difficult. Specifically, during inspection of the printed circuit board, the operator of the inspection station must determine where imperfections in the printed circuit board exist. Upon identifying such imperfections, the operator marks the imperfections and then removes the printed circuit board from the inspection station in order to rework the circuit board. Because the illumination head is positioned in close proximity to the top surface of the printed circuit board, the operator of the inspection station cannot effectuate repairs to the printed circuit board. Accordingly, the operator must mark and remove defective printed circuit boards in order to rework the same. Once the defect has been marked, rework on the printed circuit board is effectuated on another machine.\nThe present invention addresses the above-mentioned deficiencies in inspection stations by providing an illumination head for an inspection station that may be disposed 4 to 6 inches above the top surface of the printed circuit board and still generate a clear image of the inspected part. In this respect, reworking of the printed circuit board may be accomplished at the inspection station which includes the illumination head of the present invention. Accordingly, an operator using an inspection station having an illumination head of the present invention may find defects on the top surface of the printed circuit board and rework the same without removing the printed circuit board from the inspection station.\nAn illumination head for an inspection and rework station. The illumination head has an optical axis and is operative to illuminate a part. The illumination head comprises a backlight having an aperture formed therein which is coaxially aligned with the optical axis of the illumination head. Further, the illumination head comprises at least one lightline disposed adjacent to the backlight. The backlight is operative to direct light onto the part along the optical axis and the lightline is operative to direct light onto the part at an angle incident to the optical axis. Accordingly, the illumination head may be disposed a prescribed distance above the part (i.e., typically four to six inches).\nIn accordance with the present invention, the aperture of the backlight may be circularly configured and the illumination head may further comprise an annular ringlight disposed between the backlight and the lightline. The ringlight is coaxially aligned with the aperture and operative to direct light onto the part along the optical axis. The lightline typically comprises two lightlines disposed on opposite sides of the illumination head. The illumination head may further comprise an iris coaxially aligned with the optical axis and configured to selectively regulate the level of illumination directed onto the part.\nThe illumination head constructed in accordance with the present invention may further include a viewing device coaxially aligned with the optical axis and operative to generate an image of the part through the aperture of the backlight. The viewing device may be a CCD camera.\nIn the preferred embodiment, the backlight is a translucent sheet configured to transmit light. The backlight is in optical communication through at least one optical fiber with a light source operative to illuminate the backlight. Similarly, the lightline may comprise a plurality of optical fibers disposed in substantially parallel relation to one another and operative to direct light upon the part at an angle incident to the optical axis.\nIn accordance with the present invention, there is provided a method of illuminating a part with an illumination head having an optical axis, a lightline, and a backlight. The method comprises illuminating the part with light directed at an angle incident to the optical axis with the lightline and illuminating the part with light directed along the optical axis with the backlight such that the illumination head is disposed a prescribed distance above the part. The method further may include coaxially aligning the aperture with the optical axis prior to illuminating the part.\nIn accordance with the present invention, the illumination head may further comprise an annular ringlight coaxially aligned with the aperture such that the method further comprises illuminating the part with light directed along the optical axis with the ring light. Additionally, the illumination head may further comprise an iris coaxially aligned with the aperture such that the method comprises controlling an amount of illumination from the backlight and the ringlight with the iris."}
{"text": "The present disclosure generally relates to enabling devices within an industrial automation system to become aware of certain attributes pertaining to the industrial automation system or a part of the industrial automation system, in which the devices are located. More specifically, the present disclosure relates to systems and methods for industrial automation devices to analyze its received data with respect to various parts of the industrial automation system or the industrial automation system as a whole."}
{"text": "The present invention relates to turbomachinery and vanes for guiding air flow through a turbomachine. More particularly, this invention relates to a fan outlet guide vane frame comprising sectors of outlet guide vanes formed of different materials, in which some of the sectors are load-bearing and others are not.\nHigh-bypass turbofan engines are widely used for high performance aircraft that operate at subsonic speeds. As schematically represented in FIG. 1, a high-bypass turbofan engine 10 includes a large fan 12 placed at the front of the engine 10 to produce greater thrust and reduce specific fuel consumption. The fan 12 serves to compress incoming air 14, a portion of which flows into a core engine (gas turbine) 16 that includes a compressor section 18 containing low and high pressure compressor stages 18A and 18B to further compress the air, a combustion chamber 20 where fuel is mixed with the compressed air and combusted, and a turbine section 22 where a high pressure turbine 22A extracts energy from the combustion gases to drive the high pressure stages 18A of the compressor section 18 and a low pressure turbine 22B extracts energy from the combustion gases to drive the fan 12 and the low pressure stages 18B of the compressor section 18. A larger portion of the air that enters the fan 12 is bypassed to the rear of the engine 10 to generate additional engine thrust. The bypassed air passes through an annular-shaped bypass duct 24 that contains one or more rows of stator vanes, also called outlet guide vanes 28 (OGVs), located immediately aft of the fan 12 and its fan blades 26. The fan blades 26 are circumscribed by a fan casing 32, which in turn is surrounded by the fan cowling or nacelle 34 that defines the inlet duct 36 to the turbofan engine 10 as well as a fan nozzle 38 for the bypassed air exiting the bypass duct 24.\nThe outlet guide vanes 28 form part of a vane frame 40 that further includes inner and outer shrouds 42 and 44 at the radially inward and outward extents, respectively, of the guide vanes 28. A common construction is to form the vane frame 40 of segments, each comprising one or more vanes 28 connecting a pair of inner and outer shroud segments. The outer shroud 44 (formed by the assembly of the outer shroud segments) is secured to the fan casing 32, while the inner shroud 42 (formed by the assembly of the inner shroud segments) is secured to the core engine 16, and more particularly to an inner frame (not shown) of the core engine 16. The fan nacelle 34 is shown in FIG. 1 as attached to and supported by the core engine 16 through the outlet guide vanes 28. The guide vanes 28 have cambered airfoil shapes to modify the air flow through the bypass duct 24 to promote deswirling of the fan air, thus improving efficiency and reducing engine noise.\nBecause of its dual functions, the vane frame 40 is an important structural component whose design considerations include aerodynamic criteria as well as the ability to provide sufficient structural support and stiffness to the fan nacelle 34 for maintaining the shape of the inlet duct 36 and adequately transitioning static and dynamic loads to the engine core 16. For these reasons, it is important to select appropriate constructions, materials and assembly methods when manufacturing the vane frame 40 and its individual components, including the guide vanes 28 and inner and outer shrouds 42 and 44 and their connections to the fan casing 32 and core engine 16. Various materials and configurations for outlet guide vanes have been considered. Metallic materials, and particularly aluminum alloys, have been widely used. Composite materials have also been considered, as they offer the advantage of significant weight reduction. However, in order to be reliably capable of supporting the fan nacelle and transitioning fan nacelle loads, outlet guide vanes formed of composite materials have generally required complex attachment geometries and hardware, which increases weight and manufacturing and material costs."}
{"text": "In axial piston machines, at least one working piston is mounted in a longitudinally displaceable manner in a cylinder bore of a piston drum and forms a cylinder space with the cylinder bore. The cylinder space is alternately compressed and decompressed by the longitudinal movement of the working piston and, accordingly, alternately connected to a high-pressure reservoir and a low-pressure reservoir. When changing from the low-pressure reservoir connection to the high-pressure reservoir connection, pulsations occur which can result in substantial noise generation. To counteract this, so-called pre-compression volumes are used, which are formed by pre-compression spaces.\nAxial piston machines having pre-compression spaces or zones are known from, the prior art, for example from DE 197 06 114 C5. In this, a pre-compression volume or a reservoir element is integrated in a control plate or in a connection plate of the axial piston machine. The pre-compression volumes known from the prior art can additionally be controlled via valve devices.\nThe known pre-compression spaces are sealed or closed outside the housing of the axial piston machine, which requires additional installation space. With regard to automobile applications, the installation space is an increasingly important issue for axial piston machines.\nThe seal of the pre-compression spaces moreover poses a technical challenge owing to the pulsation.\nThe object of the disclosure is to reduce the installation space for creating pre-compression spaces."}
{"text": "1. Field of the Invention\nThe present invention relates to a recording apparatus which includes a process of forming electrostatic latent images on a changed recording medium by means of irradiation of laser beams, and more particularly, to a recording apparatus which is capable of recording multi-colored information on the recording medium with a plurality of laser beams.\n2. Description of the Prior Art\nA recording apparatus of the above kind includes, as shown, for instance, in FIG. 73, a drum-shaped photosensitive body 100 as the recording medium. In the periphery of the photosensitive body 100 there are arranged successively along the direction of rotation shown by the arrow, a first charger 101, a first exposure unit 102, a first developing unit 103, a second charger 104, a second exposure unit 105, a second developing unit 106, a transfer-stripping charger 107, a cleaner 110, and a discharger 109. One cycle of process is completed by electrically charging the photosensitive body 100 uniformly with the first charger 101, forming a second electrostatic latent image by the second exposure section 105, visualizing a second color by the second developing unit 106, carrying out a control processing if needed to equalize the amount of charges by the two color toners, though not shown, transferring dichromatic information onto a transfer material 108 with the transfer-stripping charger 107, cleaning with the cleaner 110 the toner that remains on the photos ensitive body 100 after transfer, and erasing the latent images with the discharger 109.\nHowever, the existing apparatus has the second developing unit 106 which is of contact development type so that even when there is formed a first toner image 103a which is brought out to be visible, for example, by the first developing unit 103 as shown in FIG. 74(A), there may occur a case in which a portion of the first toner image 103a is scraped off by the second developing unit 106 as shown in FIG. 74(B). Then, in response to the exposure condition of the second exposure section 105, second toner 106a may be piled up by the second developing unit 106 over the first toner image 103a as shown in FIG. 74(C).\nOn the other hand, when the first toner 103a that was scraped off by the second developing unit 106 is sent into the inside of the second developing unit 106 to be mixed with the second toner 106a as shown in FIG. 75, the life of the developer (consisting of a carrier and a toner) will undergo a sharp reduction.\nFurther, in the case of the dichromatic printing process in which both of the first developing unit 103 and the second developing unit 106 are operated in the normal development mode, the changes in the surface potential of the photosensitive body 100, the conditions of the toner on the photosensitive body 100, and so forth will change as illustrated in FIG. 76(A).\nNamely, due to charging by the first charger 001, the surface potential of the photosensitive body 100 is raised, and when the normal exposure is given using the first exposure section 102, only the information zone which is irradiated by the laser beam is maintained at a high potential to form an electrostatic latent image, leaving the outside of the information zone at a low potential. The electrostatic latent image is brought out to be visible using a negatively charged toner by the first developing unit 103. When the photosensitive body 100 is charged again in this state by the second charger 104, the surface potential of the photosensitive body 100 returns to nearly the level of the first charged state and the surface toner on the electrostatic latent image is transformed to a state in which it is charged positively by the appended charges.\nNext, when the photosensitive body 100 is exposed normally by the second exposure section 105, there is formed an electrostatic latent image with high potential in the information zone, and at the same time there remains the image that was visualized in the past by the first developing unit 103. Further, an electrostatic latent image is brought out visible by the second developing unit 106 in a second exposure using negatively charged toner. A small amount of the toner is attached also to the electrostatic latent image due to the first exposure\nThe electrostatic latent image that is brought out to be visible in this manner by the two normal development modes is transferred onto the transfer material 108.\nIn addition, in the case of the dischromatic printing process in which the first developing unit 103 is operated in the inverted development mode and the second developing unit 106 is operated in the normal development mode, the surface potential of the photosensitive body 100 due to charging by the first charger 101 is raised, and an inverted exposure is carried out by the first exposure section 102 as shown in FIG. 76(B), bringing the information zone alone in low potential to form an electrostatic latent image, with the area outside of the information zone maintained at high potential. The electrostatic latent image is brought out to be visible by the first developing unit 103 due to positively charged toner. When the photosensitive body 100 is charged in this state again by the second developing unit 104, the surface potential of the photosensitive body 100 returns to approximately the level of the first charging.\nNext, when the photosensitive body 100 is exposed normally by means of the second exposing section 105, the information zone becomes an electrostatic latent image with high potential, and the image that was brought out visible by first developing unit 103 remains as is. Then, the electrostatic latent image due to the second exposure is brought out visible by the second developing unit 106 with negatively charged toner, and a small amount of the toner is attached also to the electrostatic latent image due to the first exposure After carrying out a pre-transfer charging with a charger which is not shown in order to give the same polarity to the electrostatic latent images that are brought out visible in this manner by the inverted development mode and the normal development mode, each of the electrostatic latent images that are brought out to be visible is transferred onto the transfer material 108.\nIn the case of the conventional dichromatic printing process by the combination of the normal-normal development modes or the dichromatic printing process by the combination of the inverted-normal development modes, there is necessarily involved a process of charging a toner with the charge that has the polarity that is opposite to the polarity of the toner.\nIn particular, in the dichromatic printing process by the combination of the inverted-normal development modes, the polarity of the toner used varies for each development mode so that there is an inconvenience in that in order to transfer simultaneously both electrostatic latent images that are brought out to be visible onto the transfer material 108, there has to be given a pre-transfer charging to invert the polarity of one of the toners. Moreover, when the dichromatic printing process is employed in which development is carried out in the inverted mode after development in the normal mode, there also arises the necessity of carrying out a pre-transfer charging.\nFurthermore, in the dichromatic printing process of the combination of the normal-normal development modes, the toner\nis the same in each of the developing units. However, it is inevitable to have the opposite charge on the toner, at the time of recharging with the second charger 104, as shown in FIG. 76(A).\nWhen the opposite charge appears on the toner, although each image is transferred later with corona of respective polarity, it is clear that the efficiency for each is lower than that for the ordinary monochromatic transfer.\nHowever, when a high resistance is given to the toner in order to enhance the transfer efficiency and to secure a stable development in a humid atmosphere, there arises a problem that the toner which sits on the photosensitive body inverts the polarity so that it is difficult to invert the polarity even with the reversed charging.\nIn addition, when the thickness of the toner layer on the photosensitive body is large, the toner layer is laminated in multiple layers rather than in a single layer. In such a case, when the top layer in particular is inverted, it prevents the transfer of the opposite charge to the inner toner layers, so that there is a problem that the toner polarity in the lower layers is difficult.\nMoreover, the existing color copier in practice is of the type in which an image is transferred onto a transfer paper or an intermediate transfer drum for each color, and this process is repeated, to complete the full color print, so that this method can also be applied to the recording apparatus of the type under consideration.\nHowever, in that case, the copying speed will have to be reduced sharply since a sheet of copy is obtained by repeating the process similar to the above.\nFurthermore, in the existing recording apparatus of the above kind, when printing is done in only one color, if, for instance, while the apparatus is in printing operation in a first color, a second color is designated, then the printing operation in the second color will be initiated by temporarily interrupting the rotational driving of the photosensitive body simultaneous with the completion of printing operation in the first color. Therefore, the copying speed will have to be reduced in some cases."}
{"text": "Electrode catheters have been in common use in medical practice for many years. They are used to stimulate and map electrical activity in the heart and to ablate sites of aberrant electrical activity. In use, the electrode catheter is inserted into a major vein or artery, e.g., femoral artery, and then guided into the chamber of the heart which is of concern. Within the heart, the ability to control the exact position and orientation of the catheter tip is critical and largely determines how useful the catheter is.\nBidirectional catheters have been designed to be deflectable in one direction by one puller wire and in the opposite direction within the same plane by a second puller wire. In such a construction, the puller wires extend into opposing off-axis lumens within the tip section of the catheter. So that the tip section can bend in both directions in the same plane, the puller wires and their associated lumens are located along a diameter of the tip section. For ablation catheters, electrode lead wires are also provided within the distal end and typically, an additional lumen is used to contain the electrode lead wires. For example, U.S. Pat. No. 6,210,407, is directed to a bi-directional catheter comprising two puller wires and a control handle having at least two movable members longitudinally movable between first and second positions. As another example, U.S. Pat. No. 6,171,277 is directed to a bidirectional steerable catheter having a control handle that houses a generally-circular spur gear and a pair of spaced apart rack gears. Each rack gear is longitudinally movable between first and second positions, whereby proximal movement of one rack gear results in rotational movement of the spur gear, and correspondingly distal movement of the other rack gear. Also known is U.S. Pat. No. 6,198,974 which is directed to a bi-directional electrode catheter comprising a control handle. At their proximal ends, two pairs of puller wires are attached to movable pistons in the control handle. Each piston is controlled by an operator using a slidable button fixedly attached to each piston. Movement of selected buttons results in deflection of the tip section into a generally planar “U”- or “S”-shaped curve. Further known is U.S. Pat. No. 5,891,088 directed to a steering assembly with asymmetric left and right curve configurations. Proximal ends of left and right steering wires are adjustably attached to a rotatable cam housed in a control handle. The rotatable cam has first and second cam surfaces which may be configured differently from each other to accomplish asymmetric steering.\nAlso known are control handles that provide a greater degree of deflection in the catheter tip. For example, U.S. Pat. No. 7,377,906, the entire disclosure of which is hereby incorporated by reference, has increased throw capacity through the use of pulleys around which puller wire travel for minimized offset angle between the puller wire and the longitudinal axis of the control handle while maximizing the travel distance of that puller wire for any given distance traveled by the pulley drawing the puller wire. Suitable tensile puller members are described in U.S. Patent Publication No. 2008/0103520, the entire disclosure of which is also hereby incorporated by reference.\nHowever, it is desirable to provide a control handle that allows user adjustability of the maximum degree of deflection as well as deflection sensitivity of the control handle to user manipulations, as needed for different uses and applications."}
{"text": "This application claims the priority of the corresponding German patent application Ser. No. 199 59 061.3 filed Dec. 8, 1999. The disclosure of the aforesaid priority application, as well as the disclosure of each and every US and/or foreign patent and/or patent application identified in the specification of the present application, is incorporated herein by reference.\nThe invention relates to improvements in methods of and apparatus for manipulating containers, and more particularly to improvements in methods of and in apparatus for manipulating containers of the type known as trays and often utilized for temporary storage of arrays or groups of rod-shaped articles such as plain or filter cigarettes, cigars, cigarillos, filter rod sections and other rod-shaped products of the tobacco processing industry.\nExamples of rod-shaped articles which can be confined in containers (hereinafter called trays for short) of the type to which the present invention pertains are filter rod sections or filter mouthpieces which can be united with plain cigarettes, cigarillos, cigars or analogous rod-shaped tobacco-containing products to form therewith filter cigarettes, cigarillos, cigars or the like. Many types of filter mouthpieces contain a rod-like filler of acetate fibers and hardened droplets of a plasticizer which bonds portions of neighboring fibers to each other to thus establish a maze of paths for the flow of tobacco smoke from the lighted end, through the tobacco-containing part, through the filter mouthpiece and into the mouth of a smoker.\nHeretofore known and presently utilized plasticizers (such as triacetin) for acetate fibers or the like require a certain period of time to set and to thus establish reliable bonds between neighboring portions of fibers in the mouthpiece of a filter cigarette or the like. Therefore, it is desirable to ensure that filter mouthpieces of unit length or multiple unit length which issue from a filter rod making machine remain unattached to plain cigarettes for certain intervals of time which are required by the plasticizer to set, i.e., which are needed to ensure that the appearance and/or other desirable characteristics of filter-tipped smokers\"\" products are not affected during transport of filter mouthpieces into and/or during their treatment in a so-called tipping machine wherein plain cigarettes and filter mouthpieces of unit length or multiple-unit length are connected to each other by so-called uniting bands (e.g., webs of cigarette paper, imitation cork or the like) to form therewith filter cigarettes or analogous rod-shaped smokers\"\" products.\nFilter rod making machines which are of the type utilized for the making of continuous filter rods ready to be subdivided into filter rod sections of unit or multiple unit length are disclosed, for example, in U.S. Pat. No. 3,974,007 (granted Aug. 10, 1976 to Greve for xe2x80x9cMETHOD AND APPARATUS FOR THE PRODUCTION OF FILTER ROD SECTIONS OR THE LIKExe2x80x9d) and in commonly owned U.S. Pat. No. 4,412,505 (granted Nov. 1, 1983 to Haxc3xcsler et al. for xe2x80x9cAPPARATUS FOR APPLYING ATOMIZED LIQUID TO A RUNNING LAYER OF FILAMENTARY MATERIAL OR THE LIKExe2x80x9d).\nFilter tipping machines which can be utilized to unite filter rod sections with plain cigarettes are known under the name MAX-S (distributed by the assignee of the present application). Reference may also be had to commonly owned U.S. Pat. No. 5,135,008 granted Aug. 4, 1992 to Oesterling et al. for xe2x80x9cMETHOD OF AND APPARATUS FOR MAKING FILTER CIGARETTESxe2x80x9d.\nThe apparatus of the present invention can be utilized to transport filter rod sections between a filter rod making machine (such as that disclosed by Greve or by Hxc3xa4usler et al.) and a filter tipping machine (such as that disclosed by Oesterling et al.). Certain presently known transporting apparatus of such character are disclosed in commonly owned U.S. Pat. No. 5,123,798 granted January 23, 1992 to Glxc3x6ssmann et al. for xe2x80x9cAPPARATUS FOR MANIPULATING TRAYS FOR CIGARETTES AND\nTHE LIKExe2x80x9d. One presently utilized embodiment of the patented apparatus is designed to manipulate cigarettes, and more specifically to accept the surplus of cigarettes advancing in the form of a mass flow from a cigarette maker to a processing machine (such as a tipping machine of the type disclosed in the patent to Oesterling et al.) and to admit cigarettes into the mass flow when the requirements of the processing machine exceed the output of the maker. Such apparatus employs trays which are filled with cigarettes when the output of the maker exceeds the requirements of the processing machine and which are emptied into the mass flow when the need arises. The just described apparatus employs storage facilities for filled trays at one or more first levels and storage facilities for empty or emptied trays at one or more second levels. A transfer unit is provided to transport trays between various levels and, when necessary, to change the orientation of the trays.\nIt is also known to employ pneumatic conveyor systems as a means for transporting filter rod sections from a filter rod maker to a filter tipping machine (such as the aforementioned MAX-S machine). As a. rule, a pneumatic conveyor employs a receiving station with magazines which contain supplies of filter rod sections and cooperate with pneumatic senders which propel filter rod sections to the magazine of a filter tipping machine. Such apparatus are rather complex, bulky and expensive because they must be provided with special feeders for admission of filter rod sections into the senders and with special transfer units which advance filter rod sections from the senders into the magazine(s) of one or more tipping machines. On their way from the maker to the magazine of a filter tipping machine, the filter rod sections are subjected to repeated mechanical stressing which is bound or apt to affect their quality if it takes place prematurely, i.e., before the various constituents and/or ingredients of the filter rod sections are ready to withstand such stressing.\nAn object of the present invention is to provide a novel and improved method of manipulating filled and empty receptacles (such as the aforediscussed trays) for groups of rod-shaped articles of the tobacco processing industry in such a way that filled trays remain intact for periods of time which are required to optimize the quality of their contents prior to admission into the next processing station.\nAnother object of the invention is to provide a method which can be resorted to with particular advantage in connection with the manipulation of trays for rod-shaped products of the type that require a certain minimum interval of time to xe2x80x9cripenxe2x80x9d or xe2x80x9cagexe2x80x9d starting with the instant of issuance from a maker and ending with the instant of admission into a processing machine or with the instant of undergoing mechanical stresses.\nA further object of the invention is to provide a novel and improved method of manipulating filter mouthpieces for tobacco smoke between a filter rod making machine and a tipping machine wherein filter rod sections are united with plain cigarettes, cigars, cigarillos or the like to form therewith filter cigarettes, cigars or cigarillos of unit length or multiple unit length.\nAn additional object of the invention is to provide a method which renders it possible to optimally reconcile the need for adequate setting of certain hardenable substances in rod-shaped smokers\"\" products and the requirements of modern high-speed processing machines or production lines wherein the rod-shaped products undergo further treatment or treatments such as uniting with other types of rod-shaped commodities, packing in containers (such as so-called soft packs and hinged-lid packs for plain or filter cigarettes) and/or others.\nStill another object of the invention is to provide an apparatus for the practice of the above outlined method.\nA further object of the invention is to provide a compact and relatively simple but highly versatile apparatus for the manipulation of empty, filled and emptied trays for groups of rod-shaped articles, such as filter rod sections, on their way from one or more makers to one or more processing stations.\nAnother object of the invention is to provide an apparatus which can select the periods of dwell of certain types of smokers\"\" products in an economical time and space-saving manner.\nAn additional object of the present invention is to provide the apparatus with a novel and improved regulating system which renders it possible to prevent premature processing of rod-shaped articles while at the same time avoiding excessive periods of dwell of such articles prior to further processing.\nOne important feature of the present invention resides in the provision of a method of manipulating trays or analogous receptacles or containers (hereinafter called trays for short) for groups of rod-shaped articles, especially for stacks of parallel plain or filter cigarettes, cigars, cigarillos, filter rod sections or analogous smokers+ products. The improved method comprises the steps of filling successive empty trays at a filling station (e.g., at the discharge end of a filter rod making machine), transporting successive filled trays from the filling station to a predetermined storage facility (normally the first of two or more storage facilities for filled trays), monitoring the periods of dwell of filled trays at the predetermined facility (hereinafter called first facility for short), conveying filled trays from the first facility to a tray emptying or evacuating station (e.g., to the magazine of a filter tipping machine if the rod-shaped articles are filter rod sections or plain cigarettes), and regulating the timing of the conveying step as a function of the monitored periods of dwell of the filled trays at the first facility.\nAs a rule, the regulating step includes initiating the conveying of filled trays to the emptying station (either directly or through one or more additional storage facilities) when the respective monitored periods of dwell at least match a predetermined minimum period of dwell (e.g., a period of dwell which suffices to ensure that the aforediscussed plasticizer for acetate fibers of standard filter mouthpieces has set and/or that the adhesive for the overlapping seams of the tubular wrappers of filter mouthpieces can stand the treatment of filter mouthpieces in the tipping machine). The conveying step of such method preferably includes conveying filled trays from the first facility to the emptying station when the respective monitored periods of dwell at the first facility at least match the predetermined period of dwell.\nThe regulating step can include comparing the monitored periods of dwell with the predetermined minimum period of dwell and delaying the conveying step when the minimum period of dwell exceeds the monitored periods of dwell. The conveying step of such method can include advancing filled trays from the first facility to the emptying station along a first path when the respective monitored periods of dwell at least match the minimum period of dwell, and advancing filled trays along a longer second path when the minimum period of dwell exceeds the monitored periods of dwell.\nThe monitoring step can include establishing and assigning to each of the filled trays information denoting the time of entry into the first facility, and monitoring the assigned information on conveying of filled trays from the first facility to thus ascertain the periods of dwell of filled trays in the first facility.\nThe method can further comprise the steps of establishing a magazine for the delivery of filled trays received, for example, from the first facility to the emptying station within a predetermined interval of time, and the regulating step of such method can include ascertaining for each filled tray the sum of the respective monitored period of time and the predetermined interval; the conveying step of such method preferably includes advancing filled trays from the first facility directly into the magazine when the respective sums at least match the predetermined minimum period of time. Such method can further include the step of ascertaining for each filled tray the aforesaid predetermined interval of time; this further step can comprise establishing and assigning to each filled tray the information denoting the time of entry into the magazine, and monitoring the assigned information upon arrival of the respective filled tray at the emptying station to thus ascertain the predetermined interval of time.\nThe just described embodiment of the method can further comprise the step of advancing filled trays from the first facility to a second storage facility when the predetermined minimum period of dwell exceeds the respective sum. Such method can further comprise the step of advancing filled trays from a third storage facility to the magazine at the emptying station for at least some of the filled trays which are advanced from the first facility to the second facility. Still further, such method can include the steps of monitoring the periods of dwell of filled trays in the second facility, ascertaining for each filled tray at the second facility the sum total of the predetermined interval of time and the monitored periods of time at the first and second storage facilities, and conveying filled trays from the second facility to the magazine when the respective sums total at least match the predetermined maximum period of time. Each step of monitoring the period of dwell of a filled tray in the second storage facility can include establishing and assigning to each of the filled trays information denoting the time of entry into the second storage facility, and monitoring the assigned information on conveying the filled tray from the second facility to thus ascertain the period of dwell of the. filled tray at the second storage facility. Still further, such method can comprise the step of advancing filled trays from the second storage facility into a third storage facility upon ascertainment that the sum total of the predetermined interval of time and the respective monitored periods of time is less than the predetermined minimum interval or period of time.\nThe step of advancing filled trays from the second storage facility to the third storage facility can be carried out in automatic response to the ascertainment that the sum total of the predetermined interval of time and of the respective monitored periods of time is less than the aforementioned predetermined minimum period or interval of time.\nSuch method can further comprise the step of advancing filled trays from the second storage facility, along a preselected bypass route (e.g., by hand), and into the third storage facility upon completion of ascertainment that the sum total of the predetermined interval of time and the respective monitored periods of time is less than the predetermined minimum period of time.\nThe method can further comprise the steps of establishing a third storage facility for retention of each filled tray during a second interval of time, establishing for each filled tray in the second storage facility a second sum total of the second and predetermined intervals of time and the respective monitored periods of time, and conveying filled trays from the second storage facility directly into the third storage facility when the respective second sum total at least matches the predetermined minimum period of time.\nThe method can further comprise the steps of establishing a third storage facility for the retention of a filled trays during a second interval of timer establishing for each filled tray in the second storage facility a second sum total of the predetermined and second intervals of time and the respective monitored periods of time, and conveying filled trays from the second storage facility, along a time-consuming route, and into the third storage facility when the respective second sum total is less than the predetermined minimum period of time. Such method can further comprise the step of ascertaining the second interval of time for each filled tray in the third storage facility including establishing and assigning to each filled tray information denoting the time of entry of a filled tray into the third storage facility and monitoring the assigned information on conveying of the filled tray from the third storage facility to thus ascertain the second interval of time.\nAnother important feature of the invention resides in the provision of an apparatus for manipulating trays for groups of rod-shaped articles which constitute or form part of smokers\"\" products. The improved apparatus comprises a tray filling unit, a tray emptying or evacuating unit including a magazine for filled trays a source of empty trays, a predetermined (first) storage facility for filled trays, a first conveyor which includes means for feeding a succession of empty trays from the source to the filling unit, a second conveyor which includes means for feeding successive filled trays from the filling unit to the first storage facility, and means for transporting filled trays from the first facility to the emptying unit. The transporting means includes means for monitoring intervals of dwell of successive filled trays in the magazine and/or in the first storage facility, and means for transferring filled trays from the first storage facility to the emptying unit along a path which is one of a plurality of different paths and which is selected in dependency upon the lengths or durations of monitored intervals of dwell of filled trays in the first storage facility and/or in the magazine.\nThe apparatus preferably further comprises means for regulating the operation of the transferring means; such regulating means includes means for comparing the monitored intervals of dwell with a predetermined minimum period of time. Such apparatus can further comprise a second storage facility which is arranged to receive filled trays from the first facility by way of the transferring means and to deliver filled trays to the magazine by way of the transferring means when the predetermined minimum period of time exceeds the monitored intervals of dwell. Still further, such apparatus can comprise a third storage facility for filled trays; the transferring means of such apparatus is preferably arranged to deliver filled trays from the third facility to the magazine.\nThe just described embodiment of the improved apparatus can further comprise means for effecting transfer of filled trays between the second and third storage facilities. The second and third storage facilities can be said to jointly provide a depository for filled trays being supplied by the transferring means from the first storage facility to the second or to the third facility and from the second or third facility to the magazine.\nThe trays art or can be provided with identifying indicia, and the monitoring means preferably includes sensors which can identify the trays on the basis of such indicia.\nThe first storage facility is or can be disposed at a first level, and the magazine can be disposed at a different second level.\nThe improved apparatus preferably further comprises a suitable receiver for emptied trays; the source of empty trays and the receiver for emptied trays are or can be disposed at different first and second levels, and the magazine and the first storage facility are preferably disposed at different third and fourth levels The (first) difference between the first and second levels can equal or approximate the (second) difference between the third and fourth levels The transferring means of such apparatus preferably includes an elevator having first and second tray conveying receptacles which are respectively disposed at fifth and sixth levels; the difference between the fifth and sixth levels preferably equals or approximates the first and/or the second difference.\nThe improved apparatus can further comprise a second source of empty trays; such second source and the aforementioned second storage facility for filled trays can be disposed at different first and second levels, and the first storage facility and the first mentioned source of empty trays can be disposed at different third and fourth levels. The aforementioned elevator of the tray transferring means can be set up in such a way that its first and second tray conveying receptacles are respectively disposed at fifth and sixth levels which are selected in a manner such that the difference between the fifth and sixth levels equals or approximates the difference between the first and second levels and/or the difference between the third and fourth levels. Such selection of various levels renders it possible to further speed up and automate the transfer of empty trays to the filling station, the transfer of emptied trays from the emptying station, the transfer of filled trays from the filling station, and the introduction of filled trays into the first, second or third storage facility and from a selected one of such storage facilities into the magazine at the emptying or evacuating station.\nThe manner in which empty trays are filled with arrays of preferably parallel rod-shaped articles of the tobacco processing industry and in which such articles are evacuated from filled trays is well known in the relevant art and forms no part of the present invention.\nThe novel features which are considered as characteristic of the invention are set forth in particular in the appended claims. The improved apparatus itself, however, both as to its construction and the modes of assembling, installing and operating the same, together with numerous additional important and advantageous features and attributes thereof, will be best understood upon perusal of the following detailed description of certain presently preferred specific embodiments with reference to the accompanying drawing."}
{"text": "An operator of an aircraft must often maneuver the aircraft while on the ground. This may happen during ground operations such as when the aircraft is taxiing, being maneuvered to or from a hangar, or backing an aircraft away from a terminal.\nObstacles on the ground, such as structures, other vehicles and other obstacles, may lie in the path of a vehicle. These obstacles can be detected by a person using their sense of sight. However, in many cases, due to the dimensions of the aircraft (e.g., large wing sweep angles, distance from cockpit to wingtip, etc.) and the operator's limited field of view of the areas surrounding the aircraft, it can be difficult for an operator to monitor extremes of the aircraft during ground operations. As a result, the operator may fail to detect obstacles that are located in “blind spots” in proximity to the aircraft. In many cases, the operator may only detect an obstacle when it is too late to take evasive action needed to prevent a collision with an obstacle.\nCollisions with an obstacle can not only damage the aircraft, but can also put the aircraft out of service and result in flight cancellations. The costs associated with the repair and grounding of an aircraft are significant. As such, the timely detection and avoidance of obstacles that lie in the ground path of a vehicle is an important issue that needs to be addressed.\nAccordingly, it is desirable to provide methods, systems and apparatus that can reduce the likelihood of and/or prevent collisions with the detected obstacles. It would also be desirable to assist the operator with maneuvering the aircraft and to provide an operator with aided guidance while maneuvering the aircraft so that collisions with such obstacles can be avoided. It would also be desirable to provide technologies that can be used to detect obstacles on the ground and identify aircraft position with respect to the detected obstacles (e.g., proximity of it's wings and tail or other portions that the operator can not directly observe). It would also be desirable to provide the operator with an opportunity to take appropriate evasive action to prevent a collision from occurring between the aircraft and the detected obstacles. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "A computer network is a telecommunications network that allows computers to exchange data. Network devices that originate, route, and terminate the data are called network nodes. Network nodes can include hosts, such as personal computers, phones, and servers, as well as networking hardware. In computer networks, network nodes pass data to each other along data connections. Data is typically transferred in the form of packets. Connections between network nodes are established using various media, such as fiber optic cable, coaxial cable, and wireless links.\nComputer network analytics can be used to monitor the performance of a network (for example, quality of service, network congestion, and network resilience), to monitor and enforce network security, to provide visualizations of network operation, and to support network configuration activities. Analytics applications that operate on data obtained from a network typically require the network to be configured to generate the necessary input data, in addition to aggregating that data and performing queries over that data. One aspect that typical computer network analytics applications have in common is that users need to decide beforehand which information to generate and which aggregation queries to perform. Typically, a user needs to initiate a separate query or analytics task and configure data sources to generate the data that is needed."}
{"text": "PCT Publications WO2008/151149, WO2010/06032, WO2011/150410, WO2011/150411, WO2012/061647, and WO2012/106560 disclose oils and methods for producing those oils in microbes, including microalgae. These publications also describe the use of such oils to make oleochemicals and fuels.\nTempering is a process of coverting a fat into a desired polymorphic form by manipulation of the temperature of the fat or fat-containing substance, commonly used in chocolate making.\nCertain enzymes of the fatty acyl-CoA elongation pathway function to extend the length of fatty acyl-CoA molecules. Elongase-complex enzymes extend fatty acyl-CoA molecules in 2 carbon additions, for example myristoyl-CoA to palmitoyl-CoA, stearoyl-CoA to arachidyl-CoA, or oleyl-CoA to eicosanoyl-CoA, eicosanoyl-CoA to erucyl-CoA. In addition, elongase enzymes also extend acyl chain length in 2 carbon increments. KCS enzymes condense acyl-CoA molecules with two carbons from malonyl-CoA to form beta-ketoacyl-CoA. KCS and elongases may show specificity for condensing acyl substrates of particular carbon length, modification (such as hydroxylation), or degree of saturation. For example, the jojoba (Simmondsia chinensis) beta-ketoacyl-CoA synthase has been demonstrated to prefer monounsaturated and saturated C18- and C20-CoA substrates to elevate production of erucic acid in transgenic plants (Lassner et al., Plant Cell, 1996, Vol 8(2), pp 281-292), whereas specific elongase enzymes of Trypanosoma brucei show preference for elongating short and midchain saturated CoA substrates (Lee et al., Cell, 2006, Vol 126(4), pp 691-9).\nThe type II fatty acid biosynthetic pathway employs a series of reactions catalyzed by soluble proteins with intermediates shuttled between enzymes as thioesters of acyl carrier protein (ACP). By contrast, the type I fatty acid biosynthetic pathway uses a single, large multifunctional polypeptide.\nThe oleaginous, non-photosynthetic alga, Protetheca moriformis, stores copious amounts of triacylglyceride oil under conditions when the nutritional carbon supply is in excess, but cell division is inhibited due to limitation of other essential nutrients. Bulk biosynthesis of fatty acids with carbon chain lengths up to C18 occurs in the plastids; fatty acids are then exported to the endoplasmic reticulum where (if it occurs) elongation past C18 and incorporation into triacylglycerides (TAGs) is believed to occur. Lipids are stored in large cytoplasmic organelles called lipid bodies until environmental conditions change to favor growth, whereupon they are mobilized to provide energy and carbon molecules for anabolic metabolism."}
{"text": "This invention relates to boiling and condensing heat transfer type cooling devices, and more particularly to such boiling type cooling devices for cooling power semiconductor switching elements of inverters and choppers, etc., of electric components mounted on railroad vehicles.\nFIG. 1 is a sectional view of a conventional boiling and condensing heat transfer type cooling device (hereinafter referred to as a boiling type cooling device or simple as a boiling cooling device) for cooling thyristors, which is disclosed, in Mitsubishi Denki Giho (Mitsubish Electric Corporation's Techinical Journal), Vol. 48, No. 2, 1974, p. 231.\nIn FIG. 1, flat thyristor elements 1 coupled to fins 2 in pressured contact therewith are accomodated within an evaporator 3 containing a liquid fluorocarbon 4. A condenser 5 is disposed above the evaporator 3 such that the fluorocarbon vapor 4a generated within the evaporator 3 is introduced via a vapor conduit pipe 6 to the condenser 5. A liquid return conduit pipe 7 carries to the evaporator 3 the liquid fluorocarbon 4 condensed from the fluorocarbon vapor 4a within the condenser 5.\nThe operation of the device of FIG. 1 is as follows. When the thyristors 3 are in operation, the electrical power loss taking place therein results in the generation of an amount of heat as large as several hundred watts. The heat thus generated is transferred to the fluorocarbon 4 via the fins 2 in pressured contact with the cooled surfaces of the thyristors elements 1. The flux of heat flowing through the fins 2 reaches as high a value as 10.sup.5 W/m.sup.2. As a result, the fluorocarbon 4 boils, and the thyristors elements 1 are cooled by means of the so-called boiling fluorocarbon cooling. The fluorocarbon vapor 4a generated within the evaporator 3 by the boiling of the fluorocarbon proceeds via the vapor pipe 6 into the condenser 5 to be cooled and condensed therein by means of an exterior cooler fan (not shown), and then returns in the form of liquid fluorocarbon 4 to the evaporator 3 via the liquid return pipe 7. The fluorocarbon boils and condenses repeatedly as described above, such that the thyristors 1 are cooled efficiently.\nFIG. 2 shows the principle of the conventional boiling type cooling device disclosed in the above mentioned journal (Mitsubishi Denki Giho vol. 148, No. 2). The heat generating element 1 such as thyristors is immersed in the pool of cooling medium, such as fluorocarbon R113 (trifluorotrichloroethane C.sub.2 Cl.sub.3 F.sub.3) contained in the evaporator 3. When the fluorocarbon boils or evaporates, the pressure within the evaporator 3 increases. Thus, the vapor generated within the evaporator 3 proceeds into the condenser 5 to be condenser therein. The latent heat of evaporation absorbed by the vapor within the evaporator 3 and carried by it through the vapor pipe 6 is discharged in the condenser 5 by the condensing vapor. The condensed cooling medium 4 returns to the evaporator 3 via the liquid return pipe 7. Thus, continuous heat transport by means of evaporation and condensing is established, and the thyristor elements 1 are cooled continuously.\nFIG. 3a and 3b show another conventional boiling type cooler device disclosed, for example, pressed Japanese patent publication No. 59-41307. Semiconducts 1 are in into contact with the evaporator 3, and the heat generated in the semiconducts 1 is transmitted to the liquid cooling medium 4 via the cooling fins 3a disposed within the evaporator 3. The vapor 4a generated in the evaporator 3 is guided via the vapor pipe 6 to the condenser 5 having radiation pipes 5a provided with radiation fins 5b. The condensed liquid cooling medium 4 returns to the evaporator 3 via the liquid return pipe 7.\nFIG. 4 shows a deaeration device for the cooling medium (e.g., fluorocarbon 113) of a conventional boiling type cooling device. Coupled to a deaeration bath 8 containing cooling medium 4 are a plurality of pipes 11, 13, and 16, having valves 12, 14, and 17, respectively a sealed container 15 accomodating a heat-generating electrical device such as thyristor elements is coupled to the outer end of the pipe 16. Further a branch pipe 18 having a valve 19 is coupled to the cooling medium sealing pipe 16. The bath 8 is provided with an agitator 9 comprising a motor 9a, a shaft 9b, and agitation wings or spoons 9c. Further, the bath 8 is surrounded by a cooling shell (cooling box) 10.\nThe deaeration is effected as follows. First, keeping valves 12 and 17 in the closed state, valve 14 is opened to exhaust the interior of the bath 8 by means of an exhaust pump (not shown) coupled to pipe 13. Next, valve 14 is closed and valve 12 is opened to introduce liquid fluorocarbon 4 into the bath 8 via the pipe 11. After the bath 8 is filled to a predetermined level with the liquid fluorocarbon 4, valve 12 is closed again. Then, the liquid fluorocarbon 4 is cooled by the cooling shell 10 and is agitated by the agitator 9. Thereafter, valve 14 is opened to remove the non-condensing gas such as air from the liquid fluorocarbon 4. When the deaeration is over, valve 14 is closed.\nNext, valve 19 is opened to exhaust the interior of the electrical device container 15 via the branch pipe 18. Thereafter, valve 19 is closed and valve 17 opened such that the liquid fluorocarbon 4 is introduced into the electrical device container 15.\nThe above conventional boiling type cooling devices all utilize fluorocarbons as the cooling medium for boiling and condensing heat transfer. However, recent research has shown that the fluorocarbons discharged into the atmosphere reach the stratosphere and destruct the ozonosphere. Thus, there is an urgent need for a substitute for fluorocarbons. Water has good cooling characteristics when boiled and is an obvious candidate for a substitute. However, the operating temperature of chopper devices, etc., (typical electrical devices mounted on railroad vehicles) usually ranges from -20.degree. C. to 80.degree. C. Thus, water may freeze and cause failure in the cooling device if it is used as the cooling medium.\nIt is conceivable to add an antifreeze such as ethylene glycol to water. This, however, still leaves the following problems unsolved. Water boils and evaporates more easily than the glycol. Thus, the concentration of water decreases in the evaporator, while that in the condenser increases. As a result, the cooling efficiency is progressively reduced. In addition, the dilute condenser water may freeze in the condenser 5 or in the liquid return pipe 7 to cause a failure thereof at a low ambient temperature."}
{"text": "This invention relates to baseball, and more particularly to a method and playing field for conducting a baseball hitting game.\nBaseball is a game that has been enjoyed by young and old for many years. Even the most casual sports fan knows the rules of baseball and almost everyone has played the game at one time in their life.\nThe object of baseball is to hit a pitched ball into the boundaries of a playing field and safely reach base. If a player can successfully round the bases, his team scores a run. The team with the most runs at the end of the game wins.\nThe present invention involves utilizing one of the basic elements of baseball--hitting a pitched ball into the baseball playing field. Sports fans have long admired the ability of athletes to hit a pitched baseball and there have been many long debates over whether a power hitter or a placement hitter is the better athlete at hitting a baseball. It has also been said that the most difficult feat in sport is to hit a round ball with a round bat.\nMany baseball game have been developed utilizing the hitting element of baseball. For example, a game known as \"Home Run Derby\" involved two players competing against each other to see who can hit the most home runs. Each game consisted of nine innings. Each inning comprised the first player batting against a pitcher until the first player made three outs, then the second player would bat until he made three outs. An out was any ball that was swung at by the player and was not hit over the fence for a home run. Whichever player had the most home runs at the end of nine innings won the game.\nWhile \"Home Run Derby\" was a very successful game and was even a television series in the 1950's, \"Home Run Derby\" rewarded power hitters and did not give appropriate credit to batters who could not hit home runs but could control the placement of the ball in the playing field.\nIt is an object of the present invention to create a baseball hitting game that rewards equally the ability to hit a baseball a long distance and the ability to hit a baseball to a particular location. It is a further object of the present invention to create a playing field that can be used to play the game of the present invention.\nIt is a feature of the present invention that a hitter score points for hitting home runs and also scores points for hitting the ball to a particular location.\nIt is an advantage of the present invention that power hitters and placement hitters can compete on an equal basis in a contest of skill related to the ability to hit a baseball."}
{"text": "Aluminum nitride (AlN) is a refractory material melting at 2400.degree. C., which exhibits several unique chemical and physical properties, e.g., it has a density of 3.26 g/cm.sup.3, a Young's modulus of 280 GPa, a flexural strength of 400 MPa and a Knoop hardness of 1200 kg/mm.sup.2.AlN is very stable in the presence of molten metals and therefore can be used, for example, for making crucibles to hold molten metal.\nAluminum nitride is also an electrical insulator with a bandgap of 6.2 electron volts, which makes it an attractive alternative substrate material to replace alumina and beryllia in electronic packaging. The thermal expansion coefficient of AlN is nearly identical to that of silicon. This is an important property in high power applications where thermal distortion can occur between a silicon chip and the substrate due to a mismatch in the coefficients of thermal expansion of the two materials. The thermal conductivity of aluminum nitride is nearly ten times higher than alumina and approximately equal to that of beryllia. Unlike beryllia, aluminum nitride is not restricted by processing constraints because of its toxicity.\nThere is currently a great deal of interest in polymer precursor materials that can be pyrolyzed to yield ceramic materials, including aluminum nitride. Aluminum-nitrogen polymers containing no alkyl substitution on the aluminum or nitrogen atoms are described in U.S. Pat. No. 4,767,607, in which thermolysis of a mixture of aluminum chloride and hexamethyldisilazane results in formation of a polymer with the repeating unit --(Cl)Al--N(H)].sub.n. Pyrolysis of the polymer in ammonia or under vacuum yields crystalline AlN. An infusible polymeric aluminum amidimide--(NH.sub.2)Al--N(H)].sub.n that can be pyrolyzed to form AlN is described by L. Maya, Adv. Ceram. Mat., 1986, 1, 150-153.\nPolymers having the repeating unit --(R)Al--N(H)].sub.n are disclosed in U.S. Pat. No. 4,696,968 and European Patent Application 259,164. Fibers can be melt spun from the thermoplastic precursor and pyrolyzed to. AlN. L. V. Interrante et al., Inorganic Chem., 1989, 28, 252-257 and Mater, Res. Soc. Symp. Proc., 1986, 73,359-366 reported the formation of volatile crystalline precursors that can be sublimed under vacuum. A two step pyrolysis of these precursors in ammonia resulted first in an insoluble aluminum imide polymer of the form --(R)Al--N(H)].sub.n and ultimately AlN containing less than 0.5% residual carbon and oxygen. U.S. Pat. No. 4,783,430 discloses the formation of --(CH.sub.3)Al--N(H)].sub.n, which can be pyrolyzed under helium, argon or vacuum to form hexagonal AlN.\nPolymers having the repeating unit --(H)Al--N(R)].sub.n are disclosed in U.S. Pat. No. 3,505,246 and are formed by the reaction of the alane adduct H.sub.3 Al .rarw.N(C.sub.2 H.sub.5).sub.3 with a reagent such as acetonitrile. U.S. Pat. No. 4,687,657 discloses the preparation of a poly-N-alkyliminoalane that can be pyrolyzed in argon or under vacuum to form AlN.\nReacting an organic nitrile with diisobutylaluminum hydride produced organoaluminum imines having the formula RCH.dbd.NAl(i--C.sub.4 H.sub.9).sub.2, which were not isolated (L. I. Zakharkin and I. M. Khorlina, Bull. Acad. Sci. USSR, Engl. Transl., 1959, 523-524 and Proc. Acad. Sci. USSR, 1957, 112, 879). A gas containing 85% isobutene and polymers having the repeating unit ##STR1## were produced on heating the organoaluminum imine to 220.degree. to 240.degree. C. During the formation of the polymer, aluminum alkyl groups of the organoaluminum imine are eliminated as isobutene, and aluminum-nitrogen bonds are formed.\nEuropean Patent Application 331,448 discloses that AlN can be deposited on a substrate by heating the substrate and contacting it with the vapor of an aluminum-nitrogen compound having the formula CH.sub.3 (R.sup.1)Al--N(R.sup.2)(C.sub.3 H.sub.7), where R.sup.1 is alkyl and R.sup.2 is H, alkyl or aryl. A polymer of this compound is claimed, but the structure of the polymer is not disclosed."}
{"text": "Paint rollers are very popular for painting articles of all sorts, and a variety of special paint rollers have been proposed for unusual situations. U.S. Pat. No. 2,799,884 (Bedford, July, 1957) proposes a roller assembly having two roller-receiving leg portions joined at angles to each other for painting outside corners. U.S. Pat. No. 2,813,392 (McLendon, November, 1957) describes and shows a collection of three rollers arranged to embrace a portion of a cylindrical surface. U.S. Pat. No. 2,904,813 (Schleicher, September, 1959) shows a variety of arrangements of very narrow discs on supporting rods for rolling paint onto the entire peripheral surfaces of rod-like objects. Another three-roller system is proposed in U.S. Pat. No. 5,035,022 (Iuliano et al., July, 1991).\nThe paint roller arrangements of the above-mentioned patents are relatively complicated in construction, and the positions of the rollers relative to each other are such that while they may be effective to apply paint to the intended objects, applying the paint to the rollers from the usual paint tray by rolling each roller or groups of rollers in the tray is not possible. Moreover, the paint roller systems of those patents are intended primarily for painting surfaces that face the user or face laterally of the user."}
{"text": "In U.S. Pat. No. 4,530,622 to Mercer, fill is disclosed as being retained in a geotechnical structure. A plastic material mesh, which has spaced, longitudinal, oriented strands, is used to form a retainer construction for retaining fill, such as sand. Triangular compartments are formed by a number of parallel elongate portions of the mesh which are interconnected by zig-zag portions. Each zig-zag portion is contained between two respective elongate portions and is joined to an adjacent zig-zag portion at respective corners of the formed compartments. The connections are made by transversely bending strands of one portion to form loops which project out of the opposite side of the other portion, and passing a connecting member through the loops to prevent the loops from being pulled back.\nThe geotechnical structure of the Mercer patent includes a retainer construction which need not be closed on all sides and need not have a bottom or top closure. When making up the container construction, a backing of textile material may be secured against the inner side of outer faces of the geotechnical structure, depending upon the location of the structure and in fill material to be used."}
{"text": "Computers and other electronic devices contain numerous electronic components such as processors, memory and graphics products, and other integrated circuits (ICs) that give off heat. Most electronic components are heat-sensitive and may malfunction or become physically damaged if they become too hot. However, the heat threshold within which each component in a given electronic device can safely operate varies from component to component. Thus, system level cooling elements as well as cooling elements attached to individual ICs within an electronic device are vital to the functionality of many electronic devices. These cooling elements may be heat spreaders, fans, blowers, heat sinks, and others.\nSome cooling elements can be controlled manually or by a control system that is part of an electronic device. For example, a fan can be controlled to operate at varying speeds. Controllable cooling elements are advantageous in many electronic devices because they save power and reduce overall system noise by not always operating at full speed.\nSome electronic devices rely solely on system level cooling elements for their thermal management. In many electronic devices, however, system-wide cooling requires expensive and space-consuming overhead. Thus, in many instances, individual cooling solutions for some or all of the ICs within a particular electronic device are more efficient, require less space, and are less expensive than a system level cooling solution.\nMost thermal control systems that are controllable are based on the temperature of the ICs that they cool. For example, a fan's speed may be increased if a particular IC's temperature rises to an undesirable level. However, a thermal control system that is based solely on an IC's temperature is sometimes inaccurate, inefficient, and unable to recognize and react to certain trends in the IC's power usage."}
{"text": "This invention relates to a method and apparatus for applying a nanoliter quantity of liquid to a target object without solid contact.\nIn the field of testing optical fibers using an optical time domain reflectometer (OTDR), it is desirable to be able to couple a buffer fiber, which is connected to the optical I/O port of the OTDR, to the test fiber (the fiber that is to be tested). The OTDR launches pulses of optical power into the test fiber by way of the buffer fiber, and measures the level of return optical power received from the test fiber by way of the buffer fiber.\nIn order to avoid or minimize high amplitude reflections at the interface between the buffer fiber and the test fiber, it is known to provide index-matching liquid between the end faces of the two fibers. In a known machine, this is accomplished by placing the proximal (relative to the OTDR) end segment of the test fiber in a V-shaped groove formed in a fixture that is immersed in a bath of index-matching liquid. The groove extends beyond the proximal end of the test fiber so that part of the groove is not occupied by the test fiber. The distal (relative to the OTDR) end segment of the buffer fiber is placed in the groove, so that it is essentially coaxial with the test fiber end segment, and the distal end of the buffer fiber is advanced toward the proximal end of the test fiber until the two ends are very close together. The index-matching liquid then provides good optical coupling between the two fibers.\nA disadvantage of this known technique for coupling optical fibers is that it is rather cumbersome and inconvenient to carry out the coupling operation with the end segments of the fibers submerged in liquid. Further, particles of dirt can become lodged in the groove and either interfere with seating of the distal end segment of the buffer fiber or be trapped between the end faces of the two fibers and interfere with the optical coupling of the fibers."}
{"text": "The present invention relates generally to optical communications and, more particularly, to multiple symbol polarization switching differential-detection modulation formats.\nAs Internet traffic grows exponentially because of a variety of user terminals and internet services, it has prompted strong research interests on high-speed optical networks, which are the backbone infrastructure of current “Globe Village”. The data rate for optical fiber communications has moved from 10 Gbits/s to 40 Gb/s and 100 Gbits/s or even 1 Tbits/s per channel. However, one of the major challenges facing the ultra-high-data-rate dense wavelength division multiplexing (DWDM) optical fiber transmissions is the fiber nonlinearity, which causes optical signal distortions due to the various nonlinear effects in optical fiber and sets the limit of the maximal reach. DQPSK modulation is an important format for high-speed optical communications by transmitting 2 bits per symbol. At 40 Gb/s, DQPSK systems employing direct detection are attractive by having low complexity and being generally available.\nIn a digital coherent optical communication system, different types of digital signal processing (DSP) functions can be applied, to mitigate the fiber nonlinearity, such as digital back-propagation algorithms. However, the existing DSP-based fiber nonlinearity mitigation algorithms are demanding on the hardware resources, which are relatively limited and sophisticated due to the requirements of very high electronic processing bandwidth. Meanwhile, most of the existing nonlinearity mitigation algorithms show very limited system performance improvements in real experimental testing.\nIn another approach, the phase conjugation scheme has been proposed to improve the systems' nonlinearity tolerance. However, the deployment of this scheme requires at the exact middle point of the entire transmission link, thus imposing a strict and thus unpractical restrictions on the system deployment. Its spectrum efficiency would be halved because of the fiber four-wave mixing effects. In another prior effort, the polarization states for adjacent symbols are arranged in orthogonal states for improving fiber nonlinearity tolerance.\nAccordingly, there is a need for a low-cost solution to increase the nonlinearity tolerance of a direct-detection optical DQPSK system"}
{"text": "Several types of communication techniques, including telephone calls, video conference, Voice over IP (VoIP), and instant messaging are valued because they allow persons to communicate in real-time, without delays associated with communication techniques such as letter writing, email, and other types of messaging. Because people often turn to real-time communication techniques with the hope that they can establish immediate contact with persons they are trying to reach, it is often frustrating when a communication attempt of this type fails (e.g., because a person to is unavailable to respond to an incoming call or message). In such a case, the person attempting to establish the communication has few other options but to try again later, leave a voice message, leave an email, try contacting the person on another device, etc."}
{"text": "Multicomponent membranes are known as are processes for the separation of gases therewith. Such membranes employ a polymer membrane porous to gases which carries a polymer coating occluding the pores of the first. U.S. Pat. No. 4,230,463 describes in detail these multicomponent membranes and discloses a plurality of organic polymers for use as the porous separation membrane including polysulfone, copolymers of styrene and acrylonitrile, polycarbonate and cellulose acetate. Coating materials also include a variety of polymers such as polysiloxane, polyisoprene, alpha-methylstyrene and polysiloxane copolymers and polystyrene. The multicomponent membrane can be employed to separate selectively at least one gas such as hydrogen from a mixture comprising carbon monoxide, carbon dioxide, helium, nitrogen, oxygen, argon, hydrogen sulfide, nitrous oxide, ammonia and C.sub.1 to C.sub.5 hydrocarbons.\nAn earlier U.S. Pat. No. 3,350,844 provides a process for the separation of helium, hydrogen or oxygen from a mixture of gases containing nitrogen or a hydrocarbon such as methane and one of the foregoing. The process employs a polyarylene oxide membrane. The polymer film is not crosslinked and therefore has poor solvent resistance necessitating that the gas mixtures be relatively pure. Crosslinking of polymeric materials for use as semipermeable membranes is known from U.S. Pat. No. 4,353,802, however, the membranes are employed for desalination of water by reverse osmosis rather than for gas separation.\nThus, the art of which I am aware has not provided a semipermeable membrane, highly selective for certain gases in a mixture, particularly mixtures of carbon dioxide and methane. Nor, have existing membranes possessed improved strength, flexibility and solvent resistance."}
{"text": "The storage system is one of the most limiting aspects of performance of modern enterprise computing systems. Performance of hard drive based storage is determined by seek time and time for half rotation. The performance is increased by decreasing seek time and decreasing rotational latency. However, there are limits on how fast a drive may spin. The fastest contemporary drives are reaching 15,000 rpm.\nFIG. 1 illustrates a system 100 in accordance with the prior art. In the system 100, at least one computer 102-108 is coupled to a host controller 110 and 112. The host controllers 110 and 112 are coupled to a plurality of disks 114-120.\nOften, the system 100 is configured as redundant array of independent disks (RAID)-1, storing mirrored content of the disks 114-116 in the disks 118-120. The disks 114-116 are said to be mirrored by the disks 118-120.\nIncreased reliability of the computer system is achieved by duplicating the disks 114-116, the host controllers 110 and connections therebetween. Therefore, a reliable computer system is able operate at least in presence of single failure of the disks 114-120, the RAID controllers 110 and 112, the computers 102-108, and the connections therebetween. However, storage system performance may still be inadequate using the system 100. Additionally, increasing the performance of such system is currently costly and often times is not feasible.\nFurthermore, one limiting aspect of current storage systems is the fact that many types of storage devices exhibit a limited lifetime. For example, a lifetime of non-volatile memory such as flash is reduced each time it is erased and re-written. Over time and thousands of erasures and re-writes, such storage systems may become less and less reliable.\nThere is thus a need for addressing these and/or other issues associated with the prior art."}
{"text": "A. Field of the Invention\nThe present invention relates to digging implements, and in particular a backhoe.\nB. Background of the Art\nBackhoes are used extensively for excavating and for carrying objects from one area to another. Backhoes have typically been used to dig holes in the ground for trenches and for the placement of building structural components, road substructures, cables, pipes, etc.\nHeretofore, backhoes have included a first arm pivotally attached to a tractor and a second arm pivotally attached to the first arm in a scissors-like manner. A bucket is attached to the second arm for digging. Separate hydraulic actuators have typically been used to move each of the arms and the bucket. Some of these backhoes have included an extendable second arm. Furthermore, some backhoes have included a gripping device positioned opposite the bucket for gripping objects between the gripping device and the bucket. One of the gripping devices has included a gripping device statically attached to an arm of the backhoe that does not rotate relative to the arm. These gripping devices have been difficult to use because the arm and the bucket have to properly position relative to the gripping device before the gripping device can be used to pick up objects. Another gripping device includes a separate hydraulic actuator for moving only the gripping device. These backhoe are expensive to manufacture because of the cost for the extra hydraulic actuator and the cost for connecting the gripping device to the controls in the tractor. A third gripping device includes thumbs that rotate simultaneously with the bucket. These backhoes are also difficult to use because the bucket and the gripping device must be properly positioned before the gripping implement can be used. Furthermore, these backhoes are difficult to operate because the rotating gripping implement can get in the way of the rotating bucket, thus making the ground difficult to dig."}
{"text": "I. Field of the Invention\nThe present invention relates to communications. More particularly, the present invention relates to a novel and improved method and apparatus for increasing the available data rate to and from a wireless station by multiplexing signals onto multiple carriers, multiple spread spectrum code channels and/or from multiple base stations.\nII. Description of the Related Art\nThe present invention is concerned with transmitting data at rates which are higher than the capacity of a single CDMA channel. Many solutions to this problem have been proposed. One solution is to allocate multiple channels to the users and allow those users to transmit and receive data in parallel on the plurality of channels available to them. Two methods for providing multiple CDMA channels for use by a single user are described in co-pending U.S. Pat. No. 6,005,855, entitled xe2x80x9cMETHOD AND APPARATUS FOR PROVIDING VARIABLE RATE DATA IN A COMMUNICATIONS SYSTEM USING STATISTICAL MULTIPLEXINGxe2x80x9d, issued Dec. 21, 1999, and U.S. Pat. No. 5,777,999, entitled xe2x80x9cMETHOD AND APPARATUS FOR PROVIDING VARIABLE RATE DATA IN A COMMUNICATIONS SYSTEM USING NON-ORTHOGONAL OVERFLOW CHANNELSxe2x80x9d, issued Jul. 7, 1998, both of which are assigned to the assignee of the present invention and are incorporated by reference herein. In addition, frequency diversity can be obtained by transmitting data over multiple spread spectrum channels that are separated from one another in frequency. A method and apparatus for redundantly transmitting data over multiple CDMA channels is described in U.S. Pat. No. 5,166,951, entitled xe2x80x9cHIGH CAPACITY SPREAD SPECTRUM CHANNELxe2x80x9d, which is incorporated by reference herein.\nThe use of code division multiple access (CDMA) modulation techniques is one of several techniques for facilitating communications in which a large number of system users are present. Other multiple access communication system techniques, such as time division multiple access (TDMA), frequency division multiple access (FDMA) and AM modulation schemes such as amplitude companded single sideband (ACSSB) are known in the art. However, the spread spectrum modulation technique of CDMA has significant advantages over these modulation techniques for multiple access communication systems. The use of CDMA techniques in a multiple access communication system is disclosed in U.S. Pat. No. 4,901,307, entitled xe2x80x9cSPREAD SPECTRUM MULTIPLE ACCESS COMMUNICATION SYSTEM USING SATELLITE OR TERRESTRIAL REPEATERSxe2x80x9d, assigned to the assignee of the present invention and incorporated by reference herein. The use of CDMA techniques in a multiple access communication system is further disclosed in U.S. Pat. No. 5,103,459, entitled xe2x80x9cSYSTEM AND METHOD FOR GENERATING SIGNAL WAVEFORMS IN A CDMA CELLULAR TELEPHONE SYSTEMxe2x80x9d, assigned to the assignee of the present invention and incorporated by reference herein.\nCDMA by its inherent nature of being a wideband signal offers a form of frequency diversity by spreading the signal energy over a wide bandwidth. Therefore, frequency selective fading affects only a small part of the CDMA signal bandwidth. Space or path diversity is obtained by providing multiple signal paths through simultaneous links from a mobile user through two or more cell-sites. Furthermore, path diversity may be obtained by exploiting the multipath environment through spread spectrum processing by allowing a signal arriving with different propagation delays to be received and processed separately. Examples of the utilization of path diversity are illustrated in copending U.S. Pat. No. 5,101,501 entitled xe2x80x9cSOFT HANDOFF IN A CDMA CELLULAR TELEPHONE SYSTEMxe2x80x9d, and U.S. Pat. No. 5,109,390 entitled xe2x80x9cDIVERSITY RECEIVER IN A CDMA CELLULAR TELEPHONE SYSTEMxe2x80x9d, both assigned to the assignee of the present invention and incorporated by reference herein.\nThe present invention is a novel and improved method and apparatus for the transmission of high speed data in spread spectrum communication systems. In the present invention, high speed data is provided by transmitting data on multiple carrier frequencies, multiple code channels and/or from multiple base stations. In a first embodiment of the present invention, multiplexed code symbols are transmitted on a plurality of carrier frequencies from the same base station. In second embodiment, code symbols are transmitted on multiple carrier frequencies with at least one corner frequency providing the code symbols is a multiple code channels. In a third embodiment, a subset of the multiplexed code symbols are redundantly provided on a different carrier from at least one additional base station. In a fourth embodiment, multiplexed symbols as transmitted on different carriers from the same base station and are redundantly transmitted on another set of carriers from a different base station. In a fifth embodiment, code symbols are multiplexed onto carriers from a plurality of base stations for increased throughput. In a sixth embodiment, code symbols are transmitted on carriers from a first base station and redundantly provided on at least one additional base station on the same carriers as used by the first base station.\nThe present invention further describes a method wherein a receiver can demodulate data on a plurality of channels and can allocate a demodulator or set of demodulators to search for other available systems, while maintaining the current level of data throughput. In addition the present invention describes a method for receiving multiple paging channels on a plurality of frequencies and code channels."}
{"text": "1. Field of the Invention\nThe present general inventive concept relates to an image forming apparatus. More particularly, the present general inventive concept relates to a printing media loading apparatus usable with an image forming apparatus.\n2. Description of the Related Art\nGenerally, an image forming apparatus includes a printing media loading apparatus that stores a certain amount of printing media and allows a pickup unit to pick up a printing medium. The printing media loading apparatus may include a main printing media loading means, a printing medium feeding cassette, an auxiliary printing media loading means, and a multi-purpose tray that can load a sheet or less than 50 sheets of printing media.\nIn order to allow a control portion of the image forming apparatus to know whether at least one printing medium is loaded on the printing medium feeding cassette and the multi-purpose tray, a paper sensor that can detect the presence of printing medium is disposed in each of the printing medium feeding cassette and the multi-purpose tray.\nBy using the paper sensor, the control portion can detect whether there is a printing medium on the printing medium feeding cassette and the multi-purpose tray, indicate the result to a user, prevent trouble due to malfunction, and perform a certain printing on a printing medium.\nHowever, in a conventional printing media loading apparatus usable with an image forming apparatus, material cost is increased since sensors are disposed in both the printing medium feeding cassette and the multi-purpose tray to detect loading of the printing media. Also, since space is needed for disposing the two sensors, it is difficult to miniaturize the image forming apparatus."}
{"text": "This invention relates to control of operation of magnetic disk drives of the type in which the disk is divided into \"hard\" sectors; that is, one in which the format of the disk is permanently defined at the time of manufacture of the disk. In particular the disk is divided into sectors by radial index marks followed by information defining the various tracks extending circularly around the disk. Each track is thus divided by the index mark into a plurality of sectors. Within each sector there are various permanently written data used for identification of the track and of the sector and for, e.g., timing of the data read and write operations as well as areas left for the actual storage of data on the disk. In order that the head can be properly controlled to read and write the data at the appropriate time, a format sequencer is required in order to keep track of the position of the head with respect to the disk. That is, to provide a signal indicating that, for example, the head is currently juxtaposed to the data field. Clearly, it would be desirable that such a format sequencer be capable of performing as many operations with as little hardware as possible. Further, it would be desirable if such a format sequencer could be used as well as for actual format sequencing, for example, for data read and write control."}
{"text": "Many networking applications require secure and authenticated communications. SSL and its related protocols are often used to enable secure communications between a client and a server. One drawback of SSL is that the handshake required to initiate an SSL connection may require significant computing resources, slowing down client access. One solution to this problem is to offload the task of SSL processing to a network appliance, which may sit in front of a server on a network and handle SSL connection requests. The appliance may then transmit data received via the SSL communications to the server either via a nonsecure channel or via a single SSL connection using connection pooling techniques.\nHowever, this solution may not be adequate for all networks. Computing SSL handshake messages may be a processor intensive task, and thus reduce the number of appliance processor cycles available for other tasks, such as servicing existing connections, load balancing, and caching. Thus there exists a need for systems and methods which accelerate the generation and processing of SSL handshake messages on a network appliance."}
{"text": "The phytopathogenic fungus Ashbya gossypii is a filamentously growing ascomycete that was first isolated as a plant pathogen in tropical and sub-tropical regions. It infects the seed capsule of cotton plants (Ashby S. F. and Nowell W. (1926) Ann. Botany 40: 69-84) and has also been isolated from tomatoes and citrus fruits (Phaff H. J. and Starmer W. T. (1987) In \"The Yeasts\", Vol. I Rose A. H., Harrison, J. S. (eds), Academic Press, London, 123 ff; Dammer K. H. and Ravelo H. G. (1990). Arch. Phytopathol. Pflanzenschutz, Berlin 26: 71-78 Dammer and Ravelo, 1990). The infection of the seed capsule is caused by transmission of A. gossyppii mycelium pieces or spores by stinging-sucking insects and causes a disease called stigmatomycosis.\nStudies characterising the karyotype of A. gossypii have been performed (Wright, 1990; Wendland, 1993; Gaudenz, 1994, \"The small genome of the filamentous fungus Ashbya gossypii: Assessment of the karyotype\", Diploma Thesis, Department of Applied Microbiology, Biocenter, University Basel). It has been found using yeast chromosomes of precisely known length as size markers that the genome of A. gossypii has a total nuclear genome size of 8.85 Mb. Presently, A. gossypii represents the most compact eukaryotic genome, compared to genome sizes of 12.5 Mb for Saccharomyces cerevisiae (Chu et al. (1986) Science, 234:1582-1585), 31.0 Mb for Aspergillus nidulans (Brody and Carbon (1989) Proc Natl Acad Sci USA. 86:6260-6263), and 47.0 Mb for Neurospora crassa (Orbach et al.(1988) Mol Cell Biology, 8:1469-1473).\nA. gossypii is systematically grouped to the endomycetales belonging to the family of spermophthoraceae. This classification is based on the observation that the spores that develop in hyphal compartments called sporangia look like ascospores, which are defined as end products of meiosis.\nSince A. gossypii is a filamentous ascomycete, and is capable of growing only by filamentous (hyphal) growth, fungal targets found in this model organism are predictive of targets which will be found in other pathogens, the vast majority of which grow in a filamentous fashion."}
{"text": "Yo-yo players, especially beginners, have been assisted by the development of yo-yos provided with a means to automatically return the yo-yo to the player's hand before the yo-yo spins out completely. Such an arrangement is described, for example, in U.S. Pat. No. 4,332,102 to Caffrey.\nThe Caffrey patent discloses a yo-yo having a rotatable bearing pulley mounted on the axle to which the yo-yo string is attached. Adjacent the pulley section of the bearing there is provided a cylindrical friction or braking means that interacts with two clutch mechanisms. The surface of the friction or braking means has a slip resistant characteristic and is in practice one or a series of O-rings, which are subject to wear. The clutch mechanism is provided with weighting means such that when the yo-yo is thrown the clutch is released by the development of centrifugal forces. The centrifugal forces are counterbalanced by a spring-loaded force such that the clutch is activated when the yo-yo slows down. The clutch engages the cylindrical friction surface of the pulley extension while the yo-yo still has sufficient momentum to enable the automatic return of the yo-yo to the player's hand. The successful development of an automatically returning yo-yo has proven to be especially valuable to beginners. It is also viewed as a valuable assistance to less gifted or handicapped players.\nThe nature of the arrangement shown in Caffrey, however, is such that it tends to critically weaken the structural integrity of the yo-yo. The pulley bearing to gain access to the clutch mechanism housed within a yo-yo half requires part of the boss enclosing the axle to be removed in the plastic mold. To enable sufficient braking capacity to be applied, up to 80% of the plastic boss must be removed where the pulley extension friction surface meets the clutch mechanisms.\nAlso, because the pulley is also the string bearing means, problems occur when bearing lubrication applied in excess finds its way with the aid of centrifugal forces to the nearby contact area between clutch and pulley such that the clutch slips and fails to return the yo-yo successfully.\nThe arrangement disclosed in Caffrey, by linking the centrifugally operated clutch to a coaxial extension of the string bearing means, has intrinsically restricted options on varying the quality of the string bearing means. The use of a dual purpose bearing that combines a string securing means as well as a clutch interfacing means where the clutch means is operatively enclosed in the yo-yo half must by nature expand laterally along the axial member to accommodate both functions. Caffrey achieves a superior spinning automatically returning yo-yo is achieved by narrowing the string bearing means thereby reducing the area frictionally contacting the axial member. Having the centrifugally activated clutch operatively engage an integral extension of the string bearing means also necessitates the use of double-loop string attachment to inhibit the string from slipping on the bearing means and thereby reducing the clutch effectiveness. The general public have difficulty in tying a double-loop attachment.\nAlso, the Caffrey arrangement, at least in its commercial embodiments in which the clutch mechanism occupies only one yo-yo half, exhibits a weight differential between the two yo-yo halves that is believed to shorten free spinning time.\nAnother problem experienced with a conventional automatically returning yo-yo which has a static spacing between yo-yo halves is that different yo-yo manoeuvres, to be performed efficiently, require different yo-yo response tolerances. Tom Kuhn in his publication \"SB2Flight Manual\" in \"The Art of Yo-Yo Choreography\", indicates a narrower string gap is better for loop-the-loops and a wider string gap is better for complex spin tricks."}
{"text": "The present invention relates to a specimen-holder system for upright microscopes.\nUpright microscopes used at the present time have a stage upon the surface of which the specimen to be examined is placed and which is generally displaceable, for focusing purposes, in the direction of the optical axis. For this purpose, two separate adjustments, a so-called coarse adjustment and a so-called fine adjustment, are required in order, on the one hand, to adapt the focal plane to greatly varying dimensions of the specimens and, on the other hand, to permit accurate focusing in the range of high magnification.\nSince coarse and fine adjustments are effected coaxially, this focusing mechanism is relatively expensive to manufacture, particularly if value is placed on high stability of the stage. Manipulations of the specimen or displacements of the specimen upon scanning should cause the specimen to migrate as little as possible out of the focal plane of the microscope.\nIn inverse-type microscopes, the specimens to be examined are placed almost generally on a fixed stage so that the surface of the specimen in the case of reflected-light specimens always falls in the focal plane of the microscope. For small focusing strokes, merely a fine adjustment is required, whereby the lens turret can be displaced vertically.\nThis has the disadvantage that the user cannot directly observe the specimen surface which is being examined at the time, as is for instance very useful when seeking specimen regions of interest. When operating with oil immersion, the oil must be applied to the downward-facing preparation-side of the specimen slide and therefore, upon scanning of the specimen, the oil easily flows into the mechanism of preparation protection of the objective lens.\nThe focusing by displacement of the lens turret which is known in inverse-type microscopes cannot be applied directly to upright microscopes since larger focusing strokes (such as are required to enable examination of specimens of different dimensions) require the use of a special optical system. With the lenses widely used up to the present time, and designed for a finite tube length, only slight focusing strokes can be effected without negatively affecting the quality of the intermediate image.\nWest German Pat. No. 2,449,291 discloses a device for applying specimen slides to microscopes. The device consists of a fixed stage against the lower edge of which standardized transmitted-light specimen slides can be applied by means of a suction unit; as a result, specimen planes of the specimens always agree with the focal plane of the microscope used. Aside from the fact that this device requires a vacuum conduit, which is not available in the case of simple microscopes, other specimens, for instance reflected-light specimens of different dimensions, cannot be used in this apparatus. Even examinations of blood-count chambers, which are frequently viewed in clinical operations alternately with standard specimen slides, cannot be carried out on this apparatus since normal blood-count chambers differ in all dimensions (i.e., thickness, cover-glass thickness, width, length) from the corresponding values for standard specimen slides.\nWest German utility model Pat No. 1,973,676 discloses a so-called auto-levelling stage for reflected-light specimens; this auto levelling stage is placed on the actual stage of an upright microscope and within it specimens of different thickness are pressed, in each case by spring force, against the lower side of an annular diaphragm. With such auto-levelling stages, the specimen plane is moved upward by a constant amount, namely the height of the auto-levelling stage, and remains the same for all specimens which vary within the height of the auto-levelling stage. Of course, the use of this auto-levelling stage is possible only on microscopes which have a coarse adjustment or some other stage adjustment by which the focal plane of the microscope can be readjusted by that amount which corresponds to the displacement of the specimen plane which is attributable to the height of the auto-levelling stage."}
{"text": "1. Field of the Invention\nThis invention relates generally to a brush holding device, and more particularly to a brush holding device for holding a brush for making electrical contact with a commutator of a small electric motor.\n2. Description of the Prior Art\nIn a small electric motor, brushes for making electrical contact with commutators are resiliently supported by brush holding plates made of a resilient, electrically conductive material and electrically connected to the brush holding plates. In this type of brush holding device having a brush holding plate, which is constructed so as to hold the top of a brush in position by the resiliency of lanced and raised pieces, damages are often caused to the top of the brush. That is, lanced and raised pieces of an L-shaped cross section are provided as a means for holding a brush on a brush holding plate by lancing and raising a part of the brush holding plate made of an electrically conductive, resilient material, and the brush is held in position between the edges of the lanced and raised pieces by forcing the brush in between the lanced and raised pieces. This construction often causes the edges of the lanced and raised pieces to cut into the engaging surfaces of the brush which is usually of brittle nature, resulting in damages to the top of the brush. On the other hand, to reduce the contact pressure of the brush to the commutators, the recent trend is toward the use of thinner materials for the brush holding plate, and consequently toward the increased width of the brush holding plate to increase the mechanical strength of the brush holding plate itself. However, an excessive increase in the width of the brush holding plate relative to the width of the brush proper is not desirable in terms of limited space inside the motor. It is necessary, therefore, to devise an effective means for holding a brush on a brush holding plate while maintaining the mechanical strength of the brush holding plate itself."}
{"text": "Automatically calculating the volumetric breast density (VBD) from an x-ray image produced during mammography is known. The volumetric breast density is defined as the ratio of the volume of the fibroglandular tissue to the overall volume of the breast. Below, the terms “fibroglandular tissue,” “glandular tissue,” and “mammographically dense tissue” or “dense tissue” are used synonymously. On the basis of this VBD value, the breast has until now been assigned a specific breast density category using fixed thresholds, e.g., a BI-RADS value from “1” to “4” according to the classification by the American College of Radiology (ACR). By way of example, this is described in U.S. Patent Publication No. 2011/0026791 A1 and in DE 10 2006 021 042 A1.\nWomen whose breast has a high VBD value have an increased risk of getting breast cancer. The increase in this risk is partly traced back to the fact that cancerous tissue is masked by mammographically dense tissue and therefore it is not identified during mammography.\nIt is known that the masking risk does not always correlate with the VBD value. FIGS. 1 and 2 depict a breast 3 compressed between two plates 1, 2 during mammography. After passing through the tissue of the breast 3, x-ray radiation 5 emanating from an x-ray source 4 is incident on an x-ray detector 6. As depicted in FIGS. 1 and 2, the same volume of fibroglandular tissue 7, 8 may cover small masses 9 in different ways. In the example depicted in FIG. 1, the fibroglandular tissue 7, which has a specific volume, is localized at a single position such that the small mass 9 is covered by the dense tissue 7. As a result of the volume of the fibroglandular tissue 7, the depicted region of the breast 3 is characterized by a specific VBD value. By contrast, the fibroglandular tissue 8, which has an identical volume, is distributed more uniformly in the volume of the breast 3 in the example depicted in FIG. 2, and so it is less likely for the small mass 9 to be covered by the dense tissue 8. Despite the VBD values being identical, the risk of masking of the small mass 9 is lower in FIG. 2. Therefore, the sole use of the VBD value is not sufficient for an accurate description of the masking effect by mammographically dense tissue 7, 8.\nIt is for this reason that the 5th edition of the ACR BI-RADS Atlas proposes new categories “a” to “d” with a verbal description of the breast density category, defined by the visually estimated portion of fibroglandular tissue in the breast. Taking into account the masking risk was proposed for the first time in this context. Thus, the category “c” may be assigned if small masses may be obscured as a result of a heterogeneous density distribution. It is proposed that the radiologist in such a case describes the position of the dense tissue in a further sentence.\nMoreover, the nature of the breast may be described with the category “a” if the breasts are almost entirely fatty. The category “b” may be present if there are scattered areas of fibroglandular density. The category “d” may be assigned if the breasts are extremely dense, which lowers the sensitivity of mammography.\nAs depicted in FIG. 3, a glandularity map was previously calculated in act 101 using the x-ray image produced during mammography, with the glandularity denoting the proportion of the fibroglandular tissue in the overall tissue. The glandularity map defines the amount of glandular tissue in each image pixel of the x-ray image, either as a specification in millimeters or as a percentage specification, with the value lying in a range from, e.g., 1% to 50%. Then, the mean glandularity, the VBD value, is established in act 102. Subsequently, the categories “a” to “d” are determined in act 103 based purely on the VBD value. The acts are carried out either manually by a radiologist or already in an automated form with the aid of available systems. However, only experienced radiologists may undertake a reliable assessment of the masking risk and generate a complete breast density report, taking into account the masking risk, in accordance with the prescriptions of the ACR. Here, there is a risk of errors of judgment."}
{"text": "Information associated with objects is typically collected using information carrying devices affixed to the objects. Information carrying devices include bar code symbols or other optically read symbols. A laser scanner, optical imager or other device scans or images the bar code symbol and decodes information encoded within the symbol.\nOne disadvantage of bar code symbols is that they may not be altered once printed. For example, if a bar code symbol encodes information about objects contained within a box, and then one of the objects is removed from the box, the bar code symbol may not be updated. Instead, a new bar code symbol must be printed and affixed to the box. Another disadvantage of bar code symbols is that they typically must be visible to the scanner or imager. If they are obscured (e.g., within a box), the scanner/imager may not read the symbol.\nRadio frequency (RF) tags overcome these limitations of bar code symbols. Certain RF tags may be electronically rewritten with data, thereby overcoming the permanency of bar code symbols. Additionally, RF tags may be interrogated or polled through opaque surfaces, such as through boxes, to exchange data therewith.\nOne shortcoming of RF tags, however, is that they are expensive to manufacture. Often, RF tags are constructed using a small semiconductor chip with an associated antenna, both of which may be expensive to manufacture.\nU.S. Pat. Nos. 5,204,681, 5,291,205, and 5,581,257 describe radio frequency automatic identification systems that overcome some shortcomings of RF tags. These patents describe radio frequency automatic identification systems that initially detect targets having numerous radio frequency resonators, such as quartz crystals. The quartz crystals may be made by a process of heating quartz to soften it and cutting crystals of approximate size and resonant frequency. The resonators are then produced by a process where resonance is measured to determine actual resonant frequency, and then the crystals are sorted based on certain predetermined frequency windows. The resonators may then be incorporated into a variety of objects, such as in paper. Information is attributed to the target (e.g., the paper), under the RF response characteristics of the target, such as the resonant frequencies of the resonators present, and/or the spatial locations of the resonators within the target.\nIn addition to quartz crystals, thin dipoles may be employed, which may be metallizations on a plastic film substrate. Information may be attributed to a target by fabricating the target with the resonators disposed at locations to encode information under a predetermined encoding system. Readers then read the radio frequency response characteristics of a target in a field near thereto of a radiating aperture that is activated by a radio frequency source."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of electrically-driven reciprocating pumps. More particularly, the invention relates to a pump which is particularly well suited for use as a fuel pump, driven by a solenoid assembly employing a permanent magnet and a solenoid coil to produce pressure variations in a pump section and thereby to draw into and express from the pump section a fluid, such as a fuel being pumped. The invention also relates to a fuel injector assembly employing such a pump.\n2. Description of the Related Art\nA wide range of pumps have been developed for displacing fluids under pressure produced by electrical drives. For example, in certain fuel injection systems, fuel is displaced via a reciprocating pump assembly which is driven by electric current supplied from a source, typically a vehicle electrical system. In one fuel pump design of this type, a reluctance gap coil is positioned in a solenoid housing, and an armature is mounted movably within the housing and secured to a guide tube. The solenoid coil may be energized to force displacement of the armature toward the reluctance gap in a magnetic circuit defined around the solenoid coil. The guide tube moves with the armature, entering and withdrawing from a pump section. By reciprocal movement of the guide tube into and out of the pump section, fluid is drawn into the pump section and expressed from the pump section during operation.\nIn pumps of the type described above, the armature and guide tube are typically returned to their original position under the influence of one or more biasing springs. Where a fuel injection nozzle is connected to the pump, an additional biasing spring may be used to return the injection nozzle to its original position. Upon interruption of energizing current to the coil, the combination of biasing springs then forces the entire movable assembly to its original position. The cycle time of the resulting device is the sum of the time required for the pressurization stroke during energization of the solenoid coil, and the time required for returning the armature and guide to the original position for the next pressure stroke.\nWhere such pumps are employed in demanding applications, such as for supplying fuel to combustion chambers of an internal combustion engine, cycle times can be extremely rapid. Moreover, repeatability and precision in beginning and ending of pump stroke cycles can be important in optimizing the performance of the engine under varying operating conditions. While the cycle time may be reduced by providing stronger springs for returning the reciprocating assembly to the initial position, such springs have the adverse effect of opposing forces exerted on the reciprocating assembly by energization of the solenoid. Such forces must therefore be overcome by correspondingly increased forces created during energization of the solenoid. At some point, however, increased current levels required for such forces become undesirable due to the limits of the electrical components, and additional heating produced by electrical losses.\nThere is a need, therefore, for an improved technique for pumping fluids in a linearly reciprocating fluid pump. There is a particular need for an improved technique for providing rapid cycle times in fluid pumps, such as fuel pumps without substantially increasing the forces and current demands of electrical driving components."}
{"text": "Pressure swing adsorption is a well-known method for the separation of bulk gas mixtures and for the purification of gas streams containing undesirable impurities. The method has been developed and adapted for a wide range of feed gases, operating conditions, product recovery, and product purity. Most large pressure swing adsorption (PSA) systems utilize multiple parallel adsorber beds operated in staggered sequential cycles using typical process steps of feed/adsorption, pressure equalization, depressurization, evacuation, purge, and repressurization. These PSA systems are widely used in the chemical process industries for the recovery and purification of valuable gaseous products such as hydrogen, carbon oxides, synthesis gas, light hydrocarbons, and atmospheric gases.\nThe design and operation of these PSA systems can present complex engineering challenges because of the large number of variables and parameters involved. These variables and parameters may include, for example, adsorbent type, adsorbent particle size, bed length/diameter ratio, gas flow velocities, gas residence times, type of PSA operating cycle, duration of steps in the PSA cycle, number of adsorbent beds, feed gas pressure, feed gas composition, product throughput, and product purity.\nA large worldwide market exists for the supply of high-purity hydrogen in the chemical process, metals refining, and other related industries. A typical commercial method for the production of hydrogen to satisfy this market is the reforming of natural gas or other methane-rich hydrocarbon streams. The reforming process is carried out by reacting the hydrocarbon with steam and/or an oxygen-containing gas (e.g., air or oxygen-enriched air), producing a crude reformate gas containing hydrogen, carbon oxides, water, residual hydrocarbons, and nitrogen. If carbon monoxide recovery is not required and hydrogen is the main product, the carbon monoxide may be converted to additional hydrogen and carbon dioxide by the water gas shift reaction to yield a shifted synthesis gas. Hydrogen recovery from this shifted synthesis gas typically includes a multiple-bed PSA process in which each adsorbent bed uses a layer of activated carbon for the removal of CO2 and CH4 followed by a layer of zeolite or molecular sieve adsorbent for the removal CO and N2. Other hydrogen-rich gas sources that can be upgraded by PSA technology to provide a high purity hydrogen product include refinery off-gases containing hydrogen and C1-C6 hydrocarbons, and include effluent streams from hydrocarbon partial oxidation units.\nThe overall cost of hydrogen from integrated reformer/PSA systems includes both capital and operating cost components. The economic production of high-purity hydrogen requires low operating and capital costs, wherein the capital costs depend largely upon the size of the reformer and the size of the vessels containing the PSA adsorbent beds. PSA bed size typically decreases as the hydrogen productivity (i.e., the amount of hydrogen produced per unit bed volume) of the PSA system increases, and the bed size also decreases as the hydrogen bed size factor (i.e., the volume of adsorbent bed required to produce a given amount of hydrogen product) of the PSA system decreases. Clearly, a smaller bed size factor and a larger hydrogen productivity are preferred.\nHydrogen productivity and recovery in PSA systems can be increased by improved process cycles and/or improved adsorbents. For example, the use of improved rapid cycle PSA processes can improve the overall economics of hydrogen production, since the size and cost of the reformer is impacted significantly by the performance of the PSA system, and improvements in PSA hydrogen recovery result directly in a smaller reformer. Improvements in PSA hydrogen recovery also result in a reduced demand for reformer feed gas, i.e. natural gas, which constitutes the largest operating cost of the reformer.\nThere is a need in the field of hydrogen production for improved design and operating methods to reduce overall capital and operating costs. This may be achieved by the use of improved PSA systems for final hydrogen recovery and purification, particularly by the application of improved rapid cycle processes in these PSA systems. This need is addressed by the embodiments of the present invention described below and defined by the claims that follow."}
{"text": "Generally speaking, a conventional case assembly for electronic appliance relies on the connection of multiple bolts which decide the quality of the case assembly. For instance, the reliability of the case assembly depends on the connection of the bolts and the cases, however, the bolts have to be threadedly connected to the cases one by one and which is a time-consuming task. The threaded holes can be easily broken and affect the connection of the bolts which may drop from the cases. U.S. Pat. No. 7,092,520 to Fuhrmann discloses the invention related to the connection of the bolts and the cases."}
{"text": "The present invention relates to a manufacturing method of a semiconductor device and more particularly to a backside film which is formed by growing a film on both surface sides of a semiconductor substrate.\nWhen a thin film is formed on a semiconductor substrate, the film is grown on either both of the obverse surface and the reverse surface of the semiconductor substrate (double-sided growth) or only the obverse surface thereof (single-sided growth), depending on the method of growing a film, the flow of the steps in the manufacturing method, the apparatus used therein and so forth.\nFor example, in the case that a polycrystalline silicon film to fabricate a gate electrode or the like, an insulating film to form a sidewall film, an interlayer insulating film or such is to be formed, as the film is generally grown by the LP-CVD (Low Pressure-Chemical Vapour Deposition) method, the deposition of the film proceeds not only on the obverse surface side of the semiconductor substrate but also on the reverse surface side thereof.\nThe backside film of this sort is, subsequently, removed by the following reasons.\nFirstly, when another film is to be formed by the CVD method following formation of the foregoing film, the semiconductor substrate may not be able to be fixed onto a CVD apparatus satisfactorily, unless the backside film is removed. Secondly, in transporting the semiconductor substrate in the steps of a manufacturing method, the presence of any backside film may prevent the semiconductor substrate from attaching onto a transportation vehicle sufficiently. Thirdly, in the step of performing photolithography, if any backside film is left behind, focus at exposure may shift.\nFor the above reasons, when a thin film is grown on both surface sides of a semiconductor substrate in the conventional manufacturing method, the step of removing the backside film is performed before proceeding other steps. The semiconductor substrate is formed into a prescribed thickness by grinding from the reverse surface, after all other steps are completed.\nHowever, in the methods of manufacturing a semiconductor device wherein a thin film is formed through double-sided growth and its backside film part is then removed, there are occasions in which waste matter and dust are generated in the steps subsequent to the step of the removal.\nFor instance, as shown in FIG. 4(a), after element isolation regions 502 are formed on a semiconductor substrate 501, a gate oxide film 503b is grown. Following that, as shown in FIG. 4(b), a polycrystalline silicon film 503 to fabricate gate electrodes is formed into a thickness of 200 nm or so. As the LP-CVD method is generally employed for the forming method, the polycrystalline silicon film 503 is grown on both of the obverse surface side and the reverse surface side of the semiconductor substrate 501. Consequently, etching is performed to remove a part of the polycrystalline silicon film 503 formed on the reverse surface side of the semiconductor substrate 501 and thereby a structure shown in FIG. 4(c) is obtained. In this instance, a portion of the polycrystalline silicon film lying on the edge section on the obverse surface side of the semiconductor substrate 501 may be, in part, removed therewith. The gate oxide film 503b on the reverse surface side can be also removed hereat, if circumstances require.\nAfter that, as shown in FIG. 4(d), a silicide film 504 with a thickness of 200 nm or so is formed only on the obverse surface side of the semiconductor substrate, using the sputtering method.\nThe polycrystalline silicon film 503 and the silicide film 504 fabricated as described above are then etched and worked into gate electrodes, as shown in FIG. 4(c). In some cases, however, that etching conducted to form gate electrodes leaves a residue 505 behind, as shown in FIG. 4(e), which causes generation of waste matter and dust.\nFurther, after the gate electrodes are formed, for the purpose of forming sidewall films on the sidewalls of the gate electrodes, an insulating film 506 with a thickness of 250 nm or so is formed through double-sided growth by the LP-CVD method (FIG. 5(f)). Next, as shown in FIG. 5(g), the insulating film 506 lying on the side of the obverse surface is etched and sidewall films 506 are formed on the lateral faces of the gate electrodes 503a. After that, as shown in FIG. 5(h), the insulating film 506 formed on the side of the reverse surface of the semiconductor substrate is removed by means of etching, but, also on this occasion, a residue 507 may remain, as shown in FIG. 5(g), causing generation of waste matter and dust.\nIf waste matter and dust are generated, as described above, in the steps of manufacturing a semiconductor device, a sufficient yield may not be able to be attained, and besides an additional steps of etching to remove waste matter and dust may become necessary, which lowers productivity.\nFurther, because the step of removing the backside film is performed independently from the step of grinding the reverse surface of the semiconductor substrate finally, the manufacturing method may become unduly complicated and, in some cases, even satisfactory productivity cannot be achieved.\nFurther, in Japanese Patent Application Laid-open No. 266192/1997, there is disclosed a method wherein a film is formed on the obverse surface as well as the reverse surface of a wafer, and thereafter the foregoing film lying on the obverse surface of the foregoing wafer is subjected to etching, while the foregoing film lying on the reverse surface of the wafer is made to remain. Further, in the method described in that publication, as the steps of the method proceed, layers of polysilicon and other materials become overlaid on the reverse surface of the semiconductor substrate, and it is described therein that these layers are peeled off together after the step of the final heating treatment, namely, a high temperature treatment at 800xc2x0 C.-850xc2x0 C., which may exert thermal stress is completed.\nNevertheless, the main concern in that publication is thermal stress produced in the semiconductor substrate by the heat treatment when the semiconductor substrate has different numbers of films formed on the obverse surface and the reverse surface. Accordingly, this problem of thermal stress is solved by growing films equally on both of the obverse surface and the reverse surface and making qualities and thicknesses of these films on both surfaces identical, and nothing is mentioned therein to reduce the amount of waste matter and dust.\nMoreover, although the films formed on the reverse surface of the semiconductor device are peeled off together after the heating steps are completed, grinding is not described to apply to the reverse surface of the semiconductor device. In effect, in a method of the publication, even after the backside films are peeled off, another film may be grown by a method without heat application, which may lead to formation of a backside film.\nAs set forth above, waste matter and dust brought about by removing the backside film which is formed through double-sided growth have not been hitherto regarded as a serious problem. Under such circumstances, the present inventors recognized waste matter and dust of this sort can be one of prime factors to lower yield and productivity in semiconductor device fabrication. Accordingly, an object of the present invention is to suppress generation of waste matter and attain satisfactory yield and productivity. Further, another object of the present invention is to achieve an improvement in productivity by performing the step of removing the backside film, concurrently with the step of grinding the reverse surface of the semiconductor substrate.\nIn light of the above problems, the present invention provides a method of manufacturing a semiconductor device, which comprises the steps of:\nforming a first film on both of an obverse surface side and a reverse surface side of a semiconductor substrate; and\nremoving, by means of grinding, said first film formed on the reverse surface side of said semiconductor substrate as well as a part of said semiconductor substrate in depth from the reverse surface.\nMore specifically, the present invention provides a method of manufacturing a semiconductor device; which comprises the steps of\nforming a polycrystalline silicon film into a thickness of not less than 50 nm but not greater than 150 nm on both of an obverse surface side and a reverse surface side of a semiconductor substrate;\nforming, only on the obverse surface side of said semiconductor substrate, a silicide film into a thickness of not less than 50 nm but not greater than 200 nm over said polycrystalline silicon film;\nworking said polycrystalline silicon film and said silicide film into shape and thereby forming a gate electrode;\nforming, on both of the obverse surface side and the reverse surface side of said semiconductor substrate, an insulating film for sidewall formation into a thickness of not less than 50 nm but not greater than 200 nm to cover said gate electrode;\netching said insulating film for sidewall formation which is formed on the obverse surface side of said semiconductor substrate and thereby forming a sidewall film on a lateral face of said gate electrode;\nforming, only on the obverse surface side of said semiconductor substrate, an interlayer insulating film into a thickness of not less than 500 nm but not greater than 1.5 xcexcm to cover said gate electrode; and\nremoving, by means of grinding, said polycrystalline silicon film and said insulating film for sidewall formation, both of which are formed on the reverse surface side of said semiconductor substrate, as well as a part of said semiconductor substrate in depth from the reverse surface.\nIn the above manufacturing methods, any backside film formed through double-sided growth is not removed until fabrication of the obverse surface of the semiconductor substrate is completed, and, after completing the obverse surface fabrication, in the step of grinding the reverse surface of the semiconductor substrate, the backside films are also all removed by means of grinding. Further, as the method of growing a film, a method of single-sided growth, in other words, a method in which no film is allowed to grow on the reverse surface side of the semiconductor substrate, should be employed therein ,as long as circumstances permit. As a result, generation of waste matter can be well suppressed and satisfactory yield can be attained. Further, because it is unnecessary to perform any additional etching to remove waste matter, the manufacturing method becomes simplified and satisfactory productivity, provided. Furthermore, since the removal of the backside films is carried out along with final grinding of the reverse surface of the semiconductor substrate, the manufacturing method becomes simplified and productivity, enhanced.\nIn this manufacturing method of the present invention, any backside film formed through double-sided growth is not removed until fabrication of the obverse surface of the semiconductor substrate is completed, and, after completing the obverse surface fabrication, in the step of grinding the reverse surface, the backside films are removed together. As a result, even if the manufacturing method comprises the step of double-sided growth of a thin film, it can be avoided that a part of the film remains as waste matter on the wafer edge section while the obverse surface is fabricated, which facilitates to attain excellent yield and productivity."}
{"text": "1. Field of the Invention\nThe present invention related to a system and a method capable of emulating existing tape drive systems and to also remotely archive and retrieve data files via encrypted validation communication protocol.\n2. Prior Art\nIt is necessary to store and backup data for many mainframe computer installations primarily for the purpose of safekeeping critical information to be used in the event of an unexpected loss of the primary copy. The backups are often remotely stored offsite of the mainframe installation.\nAt one time, ten inch, round reel tape drives were utilized on mainframe installations. The well known tape itself consists of a thin plastic base material with a coating of ferromagnetic ferric oxide powder. The round reel tapes were physically transported to an offsite location. Periodically, the tapes would be returned and then reused.\nIn the 1980's, cartridge tape units replaced the round reel tape drives. The tape cartridge system had fewer moving parts and was less prone to failure. Additionally, the tape cartridge system occupies a smaller floor footprint and consumed less power than the round reel drives. Additionally, the media itself was improved over time. Denser recording techniques allowed the cartridges to be smaller, yet hold the same amount of data. To improve cataloging and indexing functions, and facilitate data accessibility, typically one data set is placed on one tape volume. Some tape data sets span multiple volumes while others occupy less than a single volume. This can result in a significant waste of tape as most data sets occupy only a small portion of the media and the rest of the volume remains unused. Estimates are that industry norms are for tape cartridges to be less than 50% utilized. With a cartridge tape system, the same procedures for physically pulling certain cartridges and moving them to an offsite location would be performed.\nMore recently, virtual tape servers have been introduced which place a controller between a mainframe and the cartridge tape devices and attach a disk cache area from and to which data can be read and written. The controller handles the migration of data between the disk cache and the tape media in an optimal space and time fashion. The data is actually being read from and to disks. The disks are typically faster than tape devices.\nInformation regarding tape volumes is stored in a tape catalog, maintained by a tape management system running on the host mainframe. The tape management system associates a particular tape using its primary identifier, the tape's volume serial number, with the data sets stored onto it along with its retention, or expiration date. In order to manage the re-use of tapes, the retention date indicates when the data on a tape is no longer required and at such point in time, the tape may have its data overwritten or “scratched” out. Scratch tape is a common mainframe term for a tape available to be written upon, regardless of its prior contents if any.\nA scratch list is a report that is generally prepared on a daily basis that includes all of the volume serial numbers whose retention date expired on that day. A human typically refers to this report while walking through a tape library, pulling those tapes on the report so that they may be placed into the scratch pool for reuse. The tape management system imposes a safe guard against non-expired tapes being mounted in place of a scratch tape by comparing the tape's volume serial number against its catalog expiration date. This volume serial number, in addition to being hand written onto the exterior of the tape, is on the beginning of the tape prior to the start of data set information in a section known as a “header”. When a scratch tape is mounted for writing, the tape management system inspects the tape catalog to verify that the tape is truly a scratch. If not, then it is rejected and a different scratch tape requested.\nA vault list is a report prepared at some particular time interval that includes all of the volume serial numbers that are to be removed from the tape library and physically taken offsite. Mainframe data centers have the need to move or copy data to off site locations, primarily for the for the purpose of safe keeping critical information to be used in the event of an unexpected loss of the primary copy of that information. This typically involves physical transportation of the mainframe tapes, an error prone process in that sometimes all the required tapes are not sent or sometimes a tape sent in error that is later required to be retrieved in order to complete the processing of a mainframe job. Further, the data on these tapes is typically un-encrypted and therefore vulnerable to anyone being able to read it.\nThe tape management system is primarily used to cross-reference the location of a desired data set to a tape volume serial number. It is secondarily used to manage scratch lists and vault lists.\nThe present invention advances the art by allowing its practice to be supported via an encrypted communications protocol interfacing with, and relying upon, the teachings, practices and claims disclosed in U.S. Pat. No. 6,499,108 (hereinafter synonymously referred to as “Secure Agent™” or “SA”)."}
{"text": "Plastic pipe is cut by various types of saws. Band saws make reasonably straight cuts. However, band saws are not easily transported to work sites. Hand saws are capable of cutting plastic pipe, but require substantial effort to make straight cuts that are perpendicular to the axis of the pipe. Hand saws are also slow. Powered portable saws of various types are available on job sites. These power saws are generally fast, but make cuts that are not square with the axis of the pipe.\nPlastic pipes are used in water supply systems and sewer systems. Water systems and sewage systems both employ plastic pipes with relatively thick walls. Plastic pipes with thick walls are connected together with molded couplings and fittings such as elbows and tees, and adhesives. These fittings generally have internal projections that limit the distance the end of a pipe can be inserted into the fitting. If the cut end surface of a pipe is not in a plane that is perpendicular to the pipe's center line, portions of the end of the pipe cannot be inserted as far into the fitting as required to form a strong joint without leaks. Weak joints may leak the first time they convey a fluid. There is also a chance that weak joints will hold fluids initially but fail later. Unexpected failures and leaks may cause substantial damage.\nPlastic pipes are also used for pneumatic systems. These plastic pipes generally have relatively thin walls. Plastics used to make such pipes are semi-rigid. Like the thick walled pipes, molded couplers and fittings are used. Adhesives secure the pipes to the couplers and fittings. Leaks in a pipe system for a central vacuum system, for example, will add air into the system, thereby reducing the effectiveness of the system. Such leaks in a vacuum system are difficult to locate, even if the pipes are accessible. It is therefore imperative that strong joints without leaks be formed."}
{"text": "The present invention relates to a covering for utility vehicle superstructures, and in particular coverings for lateral openings that can swing up when released. More particularly, the invention relates to utility vehicle superstructures having an opening to be covered with a tarpaulin.\nCoverings for utility vehicle superstructures, in which a tarpaulin closes off a rear wall and is designed, in particular, as an extension of the tarpaulin cloth of the roof of the utility vehicle body, are known in practice. The tarpaulin is secured laterally by cords and, if appropriate, hooks and eyelets, and, after release of the lateral fastening, can be thrown onto the roof. It is disadvantageous in these known coverings that the closing of the tarpaulin along the lateral edges is complex and, moreover, the tarpaulin can, as a rule, only be thrown up with the aid of relatively long rods or the like, in order to completely free the opening of the vehicle body.\nCoverings which automatically roll up and in which a series of spring bars in the tarpaulin are guided from the upper end of the tarpaulin to a tube provided approximately in the centre, are furthermore known in practice, the pretensioning of the springs permitting the tarpaulin to automatically coil up when the lateral connections of the tarpaulin are released, the tarpaulin, when coiled up, being brought together in the manner of a spiral. The tarpaulin can be moved down again, for example by a cord which is provided at a lower edge of the tarpaulin being pulled, the lateral and lower closure of the covering furthermore taking place via cords guided through eyelets or hooks. A disadvantage of a closed covering is that the spring forces in the tarpaulin are not eliminated, and so the covering is generally deformed outward in the region of the points at which the springs engage.\nCoverings for utility vehicle superstructures, in which a tarpaulin is provided with a weather strip on the right-hand side and left-hand side, are furthermore known in practice, the tarpaulin having to be thrown by hand onto the roof of the utility vehicle body. To open up the tarpaulin, the lateral weather strips are clamped in each case in a pivoting strip, which can be pivoted in each case about a vertical axis, and the tarpaulin is closed by means of these as well. In this case, it is possible to secure the pivoting strip by an insertable lock or the like. A disadvantage of the known covering is that the throwing-up of the tarpaulin is complex and, moreover, the transported load is secured by the closed tarpaulin only to a limited extent. In addition, the locking pins are frequently lost.\nIn order to secure a load transported in the utility vehicle body, it is known in practice to diagonally fix tightening straps to the frame of the vehicle body and to the loading platform, the fixing of the tightening straps by hand being labor-intensive and having to take place in a separate working step before the covering is closed."}
{"text": "Irritable bowel disease, Crohn's disease, and Barrett's esophagus are just some of the gastrointestinal diseases that often require biopsy or tissue samples to be taken from the gastrointestinal tract. Often, a large number of biopsy samples must be taken from various locations in the gastrointestinal tract in order to properly diagnose the disease.\nVarious current biopsy forceps, however, are designed to take only one or two samples in a single pass. Thus, during procedures that require many more tissue samples, up to as many as twenty or more samples in some cases, the forceps must be advanced into and retracted out of the gastrointestinal tract numerous times. Such advancing and retracting of the forceps is time consuming, can cause trauma to the surrounding tissue, and can create sterility issues. Accordingly, a device that minimizes the number of advancements and retractions of the forceps by acquiring and storing multiple biopsy samples in a single pass is desirable. Also desirable is a device that is relatively easy to produce and assemble."}
{"text": "1. Field of the Invention\nThe present invention is directed to a low pressure plasma generator with a localizable plasma combustion chamber.\n2. Discussion of the Background\nTreatment with low pressure plasmas is an important new method for modifying the surfaces of solid bodies. The surfaces can be, e.g., etched, i.e., partially removed, or activated, i.e., in an energy-rich state that is suitable for extensive modifications, or are coated by bonding gaseous substances. For all of these methods, the surface to be modified must be subjected to a plasma. As is well-known, a gas comprising excited molecules, radicals or ions is referred to as a plasma.\nPlasmas can be generated at low gas pressures by means of microwave radiation. A prerequisite for the formation of a plasma is an adequately high field strength of the radiation. However, the field strength is the greatest in the immediate vicinity of the source of radiation and decreases with increasing distance therefrom. Therefore, the plasma may exist only in the vicinity of the source of radiation.\nThe uniform treatment of large surfaces or surfaces with complicated shapes with a plasma causes considerable difficulties. For reasons relating to their design and their energy supply, available sources of radiation cannot be disposed at any point and at any position in a low pressure chamber. Similarly, the surface to be treated cannot be moved to specified locations in the plasma combustion chamber. Therefore, the surfaces to be treated cannot be located near the plasma source.\nThe ability to ignite and maintain a plasma at a predetermined place, where it is supposed to unfold its technological effect, is called localization. The precise localization of plasma is of great importance primarily when a large surface is to be treated uniformly. This goal can be largely reached if one can succeed in localizing a plasma linearly and moving the plasma uniformly over the surface to be treated. For this purpose, either the plasma can be localized stationarily and the substrate can be moved relative thereto or the substrate can remain stationary and the plasma is moved at right angles to its longitudinal extension. However, just the linear localization of a uniform plasma causes considerable difficulties.\nThe literature reports on various possibilities for localizing microwave plasmas. These include, among others, the ignition of the plasma behind the inlet window for the microwave (Wertheimer et al., Thin Solid Films, 115 (1984), 109), the ignition of primary transmitting aerials (Alcatel DVM, 92240 Malakoff, France, machine GIR 820), the ignition by means of local pressure differences in a vacuum chamber (IKV reports, Mr. Ludwig) and the magnetic confinement with or without the utilization of an electron cyclotron resonance absorption (EP-A 279 895). Some of these possibilities were also used for localizing large area plasmas.\nThe use of surface waveguide structures, which are mounted outside the vacuum apparatus but which are in front of a microwave permeable window, allows a large area plasma to be ignited (Kieser et al., Thin Solid Films, 118 (1984), 203).\nAll of the described methods of localization have drawbacks that stand in the way of their practical application. The drawback of the arrangement described last is that the plasma burns only directly behind the window and cannot be moved within the vacuum to any arbitrary place therein by the operator. In the case of a coating plasma, the window is also coated, a feature that can lead to an absorption and reflection of the microwaves depending on the properties of the deposited layer. Long setting-up times then become necessary owing to the repeated cleaning or exchanging of the windows.\nAn ignition at a primary transmitting aerial yields a plasma whose intensity in most cases exhibits local inhomogeneities owing to the wavelength of the transmitting frequency (e.g., with a period of 12 cm at a frequency of 2.45 GHz). In this arrangement, compensating devices such as a mechanical movement of the aerial can hardly be used owing to the design of primary transmitting aerials\nThe ability to localize a plasma by means of local pressure differences is limited to the coating of largely closed bodies. This is a suitable method for coating bottles internally. However, in trying to process a flat substrate with such a system grave technological problems arise.\nOne successful method is magnetic confinement. This process is used, e.g, in the sputter technique. However, an effective magnetic confinement in achieved only if the gyration radius of the charged particles in the plasma with respect to the free path cannot be ignored. This is the case for conventional permanent magnets made of ferrite only below pressures of about 0.1 mbar.\nWith plasmas of higher pressures--of up to a few millibars--higher etching and deposition rates can be obtained during the etching and coating process. For this reason there is a need for a method for plasma confinement that also has a good localizability at higher pressures and thus allows homogeneous etching or formation of layers at simultaneously high etching and deposition rates."}
{"text": "With the growing availability of mobile devices, individuals may be increasingly exploring the expanding world of mobile applications, games, and social networks. For example, tablet computers may be used to engage children, with and without disabilities, in learning activities through an attractive and easy to use interface and design corresponding to such tablet computers.\nTablet computers may now be equipped with a variety of assets including intuitive user interfaces such as touchscreen, wireless connectivity via multiple different protocols such as Wi-Fi and Bluetooth, image capture capabilities, position sensing and/or location determination capabilities. A variety of applications have been introduced that capitalize on the widespread acceptance of tablet computers, which may, in part, be due to the general affordability. For example, some applications may be operable to offer assistance to individuals with speech difficulties by allowing them to create verbal words and/or phrases by pressing a series of images using a tablet device. Similarly some applications may be operable to combine an interactive drawing application with a robot or other device that may provide feedback based on a user's input sequence. Accessible messages may provide assistive text typing by highlighting keyboard elements as a user types by predicting the next sequence of letters or characters.\nUnfortunately, such touch-based tools may be developed assuming that the user possesses fine motor skills and thus is capable of touching small specific regions with an appropriate level of intensity and timing. However, the assumption of fine motor skills may be unwarranted in the context of individuals having limited upper body motor control, including for example, in children with cerebral palsy (CP). For example, children with CP may sustain dysfunctions in upper extremity (UE) activities, such as reaching, grasping and manipulation. Current therapeutic interventions for UE control in children have emphasized repeated practice of functional activities in various contexts with sufficient feedback. However children with CP may have difficulty in accessing devices requiring fine motor control such as, for example, a common pinch and swipe gesture operations that may be used to interface with a tablet computer."}
{"text": "1. Field of the Invention\nThe present invention relates to an electronic calculator which is capable of operating a numerical expression in sequence of touching the keyboard thereof in accordance with the order from left to right reading along the expression to be calculated, and is capable, in case of an expression including parentheses, of visually indicating in the display unit thereof a temporary answer resultant from operating the portion of the equation between the parentheses. The present invention also relates to an electronic calculator, which visually indicates not only the temporary result derived from a part of an expression between parentheses but any temporary result obtained from an independently operable portion of an expression to be calculated, by discriminating as an arithmetic block any independently operable portion of the expression to execute in turn an arithmetic operation to that portion.\n2. Description of the Prior Art\nA conventional desktop calculator has been designed with giving importance to the simple system configuration thereof, so that some operational functions might be reduced to a certain extent. For example, when a conventional calculator is operated by touching keys 3, .times., ( , 4, +, 5, ), and = in accordance with a numerical expression 3 .times. (4 + 5) =, the calculator will display or print out only the final result 27, and not any intermediate temporary result derived from a portion of the expression, such as 9 obtained from (4 + 5).\nHowever it is often necessary for an operator to be informed of a temporary result with regard to a portion of an expression to be calculated. Otherwise, an operator must redundantly operate the calculator to obtain the temporary result. In the example described above, when an operator intends to know the result of (4 + 5) as a partial result of the expression, the operator should depress keys 4, +, 5, and = to be informed of the answer 9 visually indicated in the display unit thereof, which answer may be written down on a sheet of paper, and after that the operator will clear all previous settings in the calculator to carry out the remaining operation 3 .times. 9 = by touching keys 3, .times., 9, and =. Then the final answer 27 will be indicated in the display unit. Thus, in a conventional calculator, an operator must discretely twice operate the keyboard thereof in accordance with two numerical expressions such as 4 + 5 = and 3 .times. 9 =. In case of calculating a number of numerical expressions, an operator will be worried by obtaining a lot of the temporary results.\nCalculating a numerical expression containing an exponential term, such as 3.sup.2, an operator often intends to known the partial answer, such as 9 in the aforesaid example, resultant from the exponential operation. In that case, the first touching of keys 3, a.sup.x, 2, and = causes the answer 9 to be obtained in a conventional calculator, and then the remaining part of that expression being calculated by employing the intermediate result 9. Thus, a numerical expression containing many independently operable terms should be divided into portions to be partially calculated so as to obtain intermediate result.\nSince such a conventional calculator is capable of displaying only the final answer from a numerical expression, it becomes more difficult to check out misoperations in keying before completing the calculation of the expression including more terms. Such a conventional calculator requires relatively more careful operation in keying, causing an operator to be exhausted."}
{"text": "Advancements in media delivery systems and media-related technologies continue to increase at a rapid pace. Increasing demand for media has influenced the advances made to media-related technologies. Computer systems have increasingly become an integral part of the media-related technologies. Computer systems may be used to carry out several media-related functions. The wide-spread access to media has been accelerated by the increased use of computer networks, including the Internet and cloud networking.\nMany homes and businesses use computer networks to generate, deliver, and receive data and information between the various connected computers. Users of computer technologies continue to demand increased access to information and an increase in the efficiency of these technologies. Improving the efficiency of computer technologies is desirable to those who use and rely on computers.\nWith the wide-spread use of computers and mobile devices has come an increased presence of and continued advancements in secure data communications. For example, advancements in mobile devices allow users to make secure connections and virtual private networks (VPNs) with corporate servers. Nevertheless, benefits may be realized by providing systems and methods for improving an authentication process."}
{"text": "At present, large-sized and high-resolution display panel has become a development trend of liquid crystal display device (LCD). However, with increasing of LCD panel size, RC delay of gate wiring and data wiring also increases (as shown in FIG. 1). Excessive RC delay may cause insufficient charging or wrong charging of pixel area, which seriously affects display quality of LCD panel, such as poor display of LCD panel and abnormal image.\nIn order to solve the aforesaid problem, it is proposed in the prior art that a thickness of a metal layer can be increased. Under the condition that a capacitance is not increased, a resistance is decreased and a signal line load is reduced. However, if the metal layer is too thick, topography of array film will be affected. As a result, a film coated on the metal layer has difficulty in climbing and may break (as shown in FIG. 2), which leads to problems such as corrosion of wiring and short connection. A qualified rate of a product will be seriously reduced. At the same time, there is no planarization technology in traditional array manufacturing procedure, which limits continuous decrease of a signal line resistance."}
{"text": "1. Field of the Invention\nThis invention relates to electronic components, and particularly to connectors for liquid cooling of electronic components.\n2. Description of Background\nOne method of cooling electronic components, for example, server components disposed in, for example a server rack, is via fluid circulated through the electronic components. The fluid is typically transported to the electronic components in tubing and/or piping and connected to the component with a quick release fitting. One drawback to fluid cooling systems, however, is the risk of electrical hazard from the possible inadvertent spray or splash of fluid onto an electronic component when connecting and/or disconnecting a fitting when servicing the electronic component."}
{"text": "1. Field of the Invention\nThe present invention relates to non-asbestos friction materials of excellent performance which can be used for braking in automobiles, large trucks, railroad cars and various types of industrial equipment.\n2. Prior Art\nFriction materials used for braking in automobiles, large trucks, railroad cars and various industrial equipment are desired to have a number of performance features. These features include excellent wear resistance, a friction coefficient that is both high and stable, outstanding resistance to brake fade, no generation of undesirable noises such as squeal during brake operation, and minimal attack of the rotor serving as the counter surface.\nThe friction material in such cases is generally fabricated by mixing together primarily a fibrous base, a binder and a filler, subjecting the mixed composition to preforming, thermoforming (pressing) and, if necessary, postcuring by heat treatment, then painting, baking and polishing to give a finished article. Mixture of the starting materials may be carried out by dry mixing or by wet mixing using a solvent or water.\nHowever, during mixture of the starting materials, the fibrous substance tends to form clumps (masses of starting material), which discourages uniform dispersion. Also, because hard particles and hard fibrous substances are highly abrasive, inadequate mixture resulting in non-uniform distribution, or segregation, within the friction material leads to a poor friction performance (e.g., noise performance, fading resistance, wear resistance, resistance to counter surface attack, vibration during high-speed braking).\nTherefore, one object of the invention is to provide a non-asbestos friction material of outstanding friction performance within which a fibrous substance and hard particles can be uniformly mixed and dispersed so as to enable the capabilities of these friction material constituents to be used to their fullest advantage.\nThe inventor has learned that by incorporating within a friction material composition particles of a rubber composite composed primarily of either rubber and a fibrous substance or of rubber, a fibrous substance and hard particles, the fibrous substance or the fibrous substance and the hard particles can be blended in a uniformly dispersed state within the friction material, enabling the friction material to exhibit an excellent friction performance.\nThat is, the inventor has found that the incorporation of rubber composite particles composed primarily of at least one fibrous substance and rubber, or the incorporation of rubber composite particles composed primarily of at least one fibrous substance, at least one type of hard particle and rubber, within a non-asbestos friction material comprised of a molded and cured composition that includes a fibrous base, a binder and a filler allows the uniform mixture and dispersion within the friction material of the fibrous substance or the fibrous substance and the hard particles so that these interlock well with the other constituents to form a non-asbestos friction material in which the capabilities of each constituent are used to fullest advantage. As a result, there can be obtained high-quality non-asbestos friction materials in which counter surface (rotor) attack and the amount of wear due to segregation of the hard particles and hard fibrous substance within the friction material are reduced, noise performance is improved, a decline in the friction coefficient at high temperatures is reliably prevented, and tearing and loss of rubber during brake operation are minimized.\nIn the non-asbestos friction materials to which the present invention is directed, incorporating particles of a rubber composite composed primarily of soft rubber, a fibrous substance and, according to one embodiment, hard particles within a non-asbestos friction material composition substantially eliminates the drawbacks described above because the fibrous substance, which tends to clump and disperses with difficulty, can be uniformly mixed and dispersed, and because segregation of the hard particles and the hard fibrous substance within the friction material, which causes high abrasiveness and ultimately leads to poor friction performance, such as noise and vibrations, can be prevented.\nAlso, the rubber within the friction material forms a transfer film on the counter surface (rotor), and so the friction performance is adversely affected by too much transfer or too much abrasion. Therefore, by including in the friction material of the invention particles of a rubber composite composed primarily of either rubber and a fibrous substance or of rubber, a fibrous substance and hard particles, a balance between rubber abrasion and film transfer is achieved. This balance discourages the loss of rubber composite particles from the friction material, leaves the rubber less susceptible to tearing even when subjected to shear forces during friction, and greatly reduces the amount of wear by the friction material. Moreover, loss of these rubber composite particles is less likely to occur even at high brake temperatures, making it possible to assure good performance by the friction material even when it is subjected to friction under high-temperature, high-load, high-speed conditions.\nThus, according to a first aspect of the invention, there is provided a non-asbestos friction material comprising a molded and cured composition that includes (A) a fibrous base, (B) a binder, (C) a filler, and (D) particles of a rubber composite composed primarily of at least one fibrous substance and rubber.\nAccording to a second aspect of the invention, there is provided a non-asbestos friction material comprising a molded and cured composition that includes (A) a fibrous base, (B) a binder, (C) a filler, and (D) particles of a rubber composite composed primarily of at least one fibrous substance, at least one type of hard particle and rubber.\nThe non-asbestos friction material of the invention is made by incorporating, in a composition of primarily a fibrous base, a binder and a filler, particles of a rubber composite composed primarily of at least one fibrous substance and rubber, or of at least one fibrous substance, at least one type of hard particle and rubber.\nThe rubber used in the rubber composite particles is not subject to any particular limitation. Exemplary rubbers include acrylonitrile-butadiene rubber (NBR), styrene-butadiene rubber (SBR), chlorobutyl rubber (CBR), silicone rubber, chloroprene rubber (CR), fluororubber (FR), isoprene rubber (IR), natural rubber (NR), butadiene rubber (BR), butyl rubber (IIR), acrylic rubber (AR), urethane rubber (UR), polysulfide rubber (such as that bearing the trade name designation Polysulfide TR), ethylene-propylene rubber (EPM), ethylene-propylene-diene rubber (EPDM) and chlorosulfonated polyethylene (CSM, such as that bearing the trade name designation Hypalon). Any one or combinations of two or more thereof may be used. Of these, acrylonitrile-butadiene rubber (NBR), styrene-butadiene rubber (SBR), chlorobutyl rubber (CBR), silicone rubber, chloroprene rubber (CR) and fluororubber (FR) are preferred.\nUse of the rubber in an uncrosslinked (unvulcanized) or semi-crosslinked (semi-vulcanized) state (including partially crosslinked or vulcanized rubber) is advantageous because rubber in this state mixes well with the other constituents such as the fibrous substance and the hard particles, and because crosslinkage (vulcanization) during molding and curing helps to enhance adhesion with the other constituents. In some cases, crosslinked (vulcanized) rubber can be used.\nThe fibrous substances referred to above include both fibers and fiber components. In the practice of the invention, such substances are divided into hard fibrous substances and soft fibrous substances. Generally, fibrous substances which have the ability to abrade the counter surface (rotor) are called xe2x80x9chard,xe2x80x9d while those which lack this ability and are instead themselves transferred are called xe2x80x9csoft.xe2x80x9d Specifically, for the purposes of the present invention, hard fibrous substances are those having a Mohs hardness of preferably at least 4, more preferably at least 4.5, even more preferably at least 5, and most preferably from 5 to 10. Conversely, soft fibrous substances are those having a Mohs hardness of preferably less than 4, more preferably less than 3.5, even more preferably less than 3, and most preferably from 1 to 3.\nThe fiber length of the fibrous substance is not subject to any particular limitation, but is generally within a range of about 10 to 5,000 xcexcm, and preferably about 100 to 3,000 xcexcm. Too short a fiber length lowers the strength of the friction material, whereas too great a length increases the abrasiveness to an excessive degree, resulting in excessive wear of the counter surface. Preferably, the aspect ratio of the individual fibers (fiber length/fiber diameter) is within a range of about 10 to 2,000, and especially 10 to 1,000.\nAt least one type, and preferably 1 to 4 types, of fibrous substance is used in the practice of invention. More specifically, at least one type of hard fibrous substance or at least one type of soft fibrous substance may be used alone, although the use of at least one type of hard fibrous substance in combination with at least one type of soft fibrous substance is preferred.\nExemplary hard fibrous substances include (1) ceramic fibers, (2) natural mineral fibers, (3) glass fibers, and (4) metal fibers. Any one or combinations of two or more of these may be used. Illustrative examples of ceramic fibers include ceramic fibers composed primarily of alumina and silica, ceramic fibers composed primarily of alumina, silica and zirconia, and ceramic fibers composed primarily of silica, calcium oxide and magnesium oxide. Examples of these ceramic fibers that may be used include commercial products such as Ibiden (made by Ibiden Co., Ltd.), S-Fiber SC (made by Thermal Ceramics Division of Nippon Steel Chemical K.K.), Superwool 612 (made by Morgan Crucible) and Fiberflux (made by Toshiba Monoflux K.K.). Illustrative examples of natural mineral fibers include rock wool, wollastonite and sepiolite. Illustrative examples of metal fibers include fibers made of various types of metal, such as steel, stainless steel, bronze, copper and brass.\nExemplary soft fibrous substances include aramid fibers, carbon fibers, cellulose fibers, acrylic fibers and potassium titanate fibers. Any one or combinations of two or more of these may be used. Of these, aramid fibers such as Kevlar(copyright) made by E.I. du Pont de Nemours Co. and Twaron(copyright) made by Akzo Nobel are preferred.\nThe hard particles used in the invention may be any particles hard enough to abrade cast iron, although particles having a Mohs hardness of at least 4, preferably at least 4.5, more preferably at least 5, and most preferably from 5 to 10, are especially advantageous.\nExemplary hard particles of this type include ceramic particles, metal oxide particles and various types of nitride particles. Any one or combinations of two or more of these may be used.\nIllustrative examples include silicon carbide, zirconium oxide, zirconium silicate, alumina, silica, magnesium oxide, iron oxide, titanium oxide, silicon nitride, zinc oxide, aluminum borate and titanium diboride. Of these, silicon carbide, zirconium oxide, zirconium silicate, alumina and silica are preferred.\nThe size of the hard particles varies according to such factors as the type, hardness and shape of the particles. In general, the particles are limited to a smaller size as the hardness becomes greater, whereas a lower hardness (greater softness) allows particles of a larger size to be used. A higher hardness (larger particle size) results in greater abrasiveness, which restricts the practical range of use of the friction material. Moreover, when the particles have a sharp, angular shape in the manner of broken fragments, for example, they tend to be more abrasive then when they are spherical or surface-treated.\nThe desired effects of the invention thus tend to become smaller as the hard particles of the invention, assuming the shape to be constant, have a lower Mohs hardness and a smaller size. For example, at a Mohs hardness of 9, it is preferable to use particles having a size of 0.1 to 10 xcexcm, and especially 0.5 to 5 xcexcm. At a Mohs hardness of 5, a particle size of 1 to 500 xcexcm, and especially 5 to 200 xcexcm, is preferred.\nWhen the rubber composite particles in the invention are composed primarily of rubber and a fibrous substance, the amount of the rubber component is preferably 3 to 70% by weight, more preferably 5 to 50% by weight, and most preferably 5 to 35% by weight; and the amount of the fibrous substance is preferably 30 to 97% by weight, more preferably 50 to 95% by weight, and most preferably 65 to 95% by weight. The fibrous substance used in this case may be a hard fibrous substance, a soft fibrous substance, or a combination of both. When a combination of both a soft fibrous substance and a hard fibrous substance is used, any suitable mixing proportions may be selected. The hard fibrous substance accounts for preferably 30 to 97% by weight, and especially 50 to 95% by weight, of the fibrous substance as a whole. The soft fibrous substance accounts for preferably 3 to 70% by weight, and especially 5 to 50% by weight, of the fibrous substance as a whole.\nWhen the rubber composite particles in the invention are composed primarily of rubber, a fibrous substance and hard particles, the amount of the rubber component is preferably 3 to 70% by weight, more preferably 5 to 50% by weight, and most preferably 5 to 35% by weight; the amount of the fibrous substance is preferably 3 to 96% by weight, more preferably 50 to 94% by weight, and most preferably 65 to 94% by weight; and the amount of hard particles is preferably 1 to 50% by weight, more preferably 1 to 30% by weight, and most preferably 1 to 20% by weight. The fibrous substance used in this case may be a hard fibrous substance, a soft fibrous substance, or a combination of both. When a combination of both a hard fibrous substance and a soft fibrous substance is used, any suitable mixing proportions may be selected. The hard fibrous substance accounts for preferably 30 to 97% by weight, and especially 50 to 95% by weight, of the fibrous substance as a whole. The soft fibrous substance accounts for preferably 3 to 70% by weight, and especially 5 to 50% by weight, of the fibrous substance as a whole.\nIn addition to the foregoing fibrous substance and hard particles, the rubber composite particles of the invention may have added thereto any suitable ingredients commonly used in friction materials. Illustrative examples include, without limitation, cashew dust, rubber powder, resins such as phenol resin, metal powders such as copper, zinc or aluminum, and also mica, vermiculite, graphite, coke, molybdenum disulfide, antimony trisulfide, antimony trioxide, phosphorus-based lubricants, barium sulfate and calcium hydroxide. Of these, the addition of a layered substance such as graphite, coke, mica or vermiculite is advantageous for enhancing the mixing properties.\nThe rubber composite particles of the invention are produced by first measuring out predetermined amounts of the above-described mixture of rubber and a fibrous substance (such as a combination of a hard fibrous substance and a soft fibrous substance) or mixture of rubber, a fibrous substance (such as a combination of a hard fibrous substance and a soft fibrous substance) and hard particles, and also predetermined amounts of optional materials, if necessary. These materials are loaded into a mixer, then mixed under an applied pressure. Preferably, the starting materials are divided into about two to ten portions and added to the mixer one portion at a time, although the entire amount of the raw materials may be added all at once. Preferable mixing conditions under applied pressure are 20 to 200xc2x0 C. and 1 to 100 kg/cm2 for 1 to 30 minutes, and especially 20 to 100xc2x0 C. and 1 to 100 kg/cm2 for 1 to 30 minutes.\nNext, after mixing under applied pressure, it is advantageous for the pressure to be released to 20 kg/cm2 or less, preferably 10 kg/cm2 or less, and especially 5 kg/cm2 or less. A mixing and milling step is then carried out, giving porous (sponge-like) rubber composite particles having a bulk density of preferably not more than 1/2, more preferably from 1/2 to 1/50, and most preferably from 1/5 to 1/20, the theoretical density or specific gravity. If the amount of rubber component in the rubber composite particles is large, it may be transferred to a separate mill, such as a cutting mill, ball mill, Turbomill or jet mill, and there subjected to size reduction.\nDuring mixture of the rubber composite particles, a solvent and a binder may be added to improve adhesion of the rubber with the fibrous substance and the hard particles. The solvent may be any which is capable of swelling and dissolving the rubber. A typical example of a solvent that may be used for this purpose is o-xylene. Illustrative examples of the binder, which may be in either a liquid (solvent) or powder form, include phenol resin, epoxy resin, polyimide resin, polyamide resin and cashew oil. When a solvent is used, drying is necessary. In such cases, it is preferable for milling to be preceded by drying because this facilitates milling and provides better efficiency.\nBeing endowed with a low bulk density and having a fluffy, porous (sponge-like) nature, the rubber composite particles of the invention, when included in a friction material, are able to ensure that the friction material has the desired porosity and elasticity. In addition, they enhance other properties such as noise performance and resistance to brake fade. The rubber composite particles used in the invention have an average particle size within a range of preferably 50 to 10,000 xcexcm, more preferably 100 to 5,000 xcexcm, even more preferably 300 to 5,000 xcexcm, and most preferably 500 to 2,000 xcexcm. Rubber composite particles with too small an average size may fail to fully achieve the desired effects of the invention, whereas too large a particle size may result in poor mixture within the friction material composition.\nWhen the rubber composite material particles of the invention are mixed into the friction material composition, addition and mixture of the particles into the composition is preferably carried out with the rubber within the rubber composite particles in an uncrosslinked (unvulcanized) or semi-crosslinked (semi-vulcanized) state and together with a vulcanizing agent (crosslinking agent) such as sulfur. Subsequently molding and postcuring (heat-treating) the friction material composition makes it possible to crosslink (vulcanize) the rubber composite particles, resulting in closer adhesion with the other constituents. The vulcanizing agent is generally added in an amount of from 0.05 to 20% by weight, based on the rubber composite particles, although the addition of a vulcanizing agent (crosslinking agent) may be omitted where use is made of a self-crosslinking (self-vulcanizing) rubber. In some cases, the rubber composite particles may first be crosslinked (vulcanized), then added and into the friction material composition.\nThe rubber composite particles are added in an amount within a range of preferably 1 to 40% by weight, more preferably 1 to 25% by weight, and most preferably 10 to 25% by weight, of the overall friction material composition. The addition of too small an amount of rubber composite particles may make it impossible to achieve the desired effects of the invention; namely, to enable the uniform dispersion and mixing of the fibrous substance without clumping, to enable the uniform dispersion and mixing of the hard fibrous substance and the hard particles within the friction material without segregation, and to use to full advantage the capabilities of the respective ingredients. On the other hand, too large an amount may result in an excessive proportion of organic substances in the friction material, resulting in decreased resistance to brake fade; that is, inferior braking performance at high brake temperatures.\nIn the non-asbestos friction material of the invention, by adding particles of a rubber composite composed of rubber and a fibrous substance or of rubber, a fibrous material and hard particles to a friction material composition, fibrous substances and hard particles which have hitherto been difficult to disperse and admix, especially hard particles and hard fibrous substances having a high abrasiveness, can be uniformly mixed and dispersed into the friction material, thereby making it possible to use to the fullest advantage the capabilities of each of these constituents.\nAs explained above, the non-asbestos friction material of the invention is made by molding and curing a friction material composition which includes a fibrous base, a binder, a filler, and the above-described rubber composite particles. No particular limitation is imposed on the constituents other than the rubber composite particles, although a fibrous base (A), a binder (B), and a filler (C) such as are used in conventional friction materials may be employed\nThe fibrous base serving as component (A) may be, for example, any inorganic or organic fiber commonly used in friction materials, other than asbestos. Illustrative examples include inorganic fibers such as metal fibers (e.g., iron, copper, brass, bronze, aluminum), ceramic fibers, potassium titanate fibers, glass fibers, carbon fibers, rock wool, wollastonite, sepiolite, attapulgite and artificial mineral fibers; and organic fibers such as aramid fibers, polyimide fibers, polyamide fibers, phenolic fibers, cellulose and acrylic fibers. Any one or combinations of two or more thereof may be used.\nThe fibrous base serving as component (A) may be in the form of staple fibers or a powder. The amount of addition is preferably within a range of 5 to 89% by weight, and especially 20 to 70% by weight, based on the overall friction material composition.\nThe binder serving as component (B) may be a known binder commonly used in friction materials. Illustrative examples include phenolic resins, melamine resins, epoxy resins, various rubber-modified phenolic resins, and NBR. Any one of these or combinations of two or more thereof may be used.\nThe amount of binder (B) added is within a range of preferably 5 to 50% by weight, and especially 10 to 25% by weight, based on the overall friction material composition.\nThe filler serving as component (C) may be any organic or inorganic filler which is known to be used in ordinary frictional materials. Illustrative examples include molybdenum disulfide, antimony trisulfide, calcium carbonate, barium sulfate, magnesium oxide, cashew dust, graphite, calcium hydroxide, calcium fluoride, talc, molybdenum trioxide, antimony trioxide, zirconium silicate, iron oxide, mica, iron sulfide, zirconium oxide, metal powders, quartz, silicon dioxide, rubber powder, alumina, chromium oxide and vermiculite. Any one or combinations of two or more of these may be used.\nThe amount of filler (C) added is within a range of preferably 5 to 60% by weight, and especially 10 to 40% by weight, based on the overall friction material composition.\nThe method of making the non-asbestos friction material of the invention involves first uniformly blending components (A), (B) and (C) and the rubber composite particles in a suitable mixer such as a Henschel mixer, Loedige mixer or Eirich mixer so as to give a molding powder, and preforming the powder in a mold. The preform is then molded at a temperature of 130 to 200xc2x0 C. and a pressure of 100 to 1,000 kg/cm2 for a period of 2 to 10 minutes.\nThe resulting molded article is postcured by heat-treating at 140 to 250xc2x0 C. for 2 to 48 hours, then spray-painted, baked and polished, giving the finished article.\nIn the case of automotive disk pads, for example, production may be carried out by placing the preform on an iron or aluminum plate that has been pre-washed, surface-treated and coated with an adhesive, then shaping the preform in this state within a mold, and subsequently heat-treating, spray-painting, baking and polishing. This gives a finished disk pad.\nThe non-asbestos friction materials of the invention are highly suitable for a variety of applications, including brake linings, clutch facings, disk pads, paper clutch facings and brake shoes in automobiles, large trucks, railroad cars and various types of industrial equipment."}
{"text": "The present invention relates to a safety controller for failsafe control of safety-critical processes, and more particularly to a safety controller for failsafe disconnection of a machine or machine system. The invention also relates to a method for loading or transferring a new operating program onto such a safety controller.\nFor the purposes of the present invention, a safety controller is a device or combination of devices connected to each another, which receive process signals from sensors of a machine or machine system and which use these signals to produce output signals by means of logic operations and, if appropriate, by means of further signal or data processing steps. The output signals are supplied as control signals to actuators which carry out specific actions or reactions in the machine or machine system. One preferred field of application for safety controllers is in the field of machine safety, namely monitoring of emergency stop buttons, two-hand controllers, guard doors, light curtains, stationary or rotary condition monitors and the like. Sensors such as these are used, for example, in order to safeguard a machine which otherwise might cause a hazard to the operator. When the guard door is opened or when the emergency stop button is operated, a process signal is produced which is supplied as an input signal to the safety controller. In response to the input signal, the safety controller switches off the dangerous part of the machine in a failsafe manner, by means of a connected actuator.\nA characteristic feature of a safety controller in contrast to a “normal” controller is that the safety controller must always ensure that the process (such as the dangerous machine) being controlled is in a safe state. This requirement even applies when a malfunction occurs within the safety controller or in a device which is connected to it. Safety controllers are therefore subject to extremely stringent requirements for their own failsafety, which results in considerable additional effort during development and manufacture. Generally, safety controllers require special licensing from responsible supervisory authorities before they are used, such as, in Germany for example, from the professional societies dealing with work safety or from a technical supervisory association. The safety controller must comply with specific safety standards which are defined, such as defined in European Standard EN 954-1. The present invention takes account of these special requirements. The expression “safety controller” therefore in this case relates only to a device or a combination of devices which are approved for the control of machines, machine systems and the like in accordance with at least Category 3 of the above-mentioned European Standard.\nA programmable safety controller offers the user the capability of individually defining, in accordance with his requirements, the logical operations on the input signals with the aid of software, specifically the so-called user program. A programmable safety controller thus replaces the previously normal wiring to the individual sensors with the aid of logic switching elements. In order to make it possible to carry out this function, a programmable safety controller has an operating program which is separate from the user program and which defines the basic functional scope of the safety controller. In particular, the operating program contains program code by means of which the hardware components of the safety controller are addressed directly and are thus “brought to life”.\nFurthermore, safety control rules are also generally implemented in the operating program, which the user program calls up as prepared functional modules and which the user can configure by means of input and output signals at any given time. For example, prepared functional modules for failsafe evaluation of a two-channel emergency-stop button or of a two-channel guard door might be contained in the operating program. In the user program, the user can now only define how the provided modules, i.e. the emergency stop button and the guard door, should be logically linked to one another.\nFor safety reasons, the user has no access to the operating program, i.e. he can neither replace nor modify the operating program. In the art, the operating program is often referred to as firmware.\nWO 98/44399 discloses a method for programming a safety controller in which the safety control rules are stored in the safety controller in the form of functional modules. By means of his user program, the user can select the functional modules, he can configure them and he can logically link them to one another. This is done by means of a programmer, with which the commands for selection, configuration and logic operations on the functional modules are transferred to the safety controller. As explained above, however, it is impossible for the user to access the safety control rules implemented in the functional modules, i.e. he can neither replace them nor modify them.\nThe inhibited access to the operating program corresponds to the well accepted practice for safety controllers, since the operating program in conjunction with the hardware of the safety controller are subject to licensing by the responsible supervisory authorities. If it were possible for the user to access the combination of hardware and operating program, the manufacturer of the safety controller would not be able to guarantee the failsafety in accordance with the verified certification, according to the general opinion.\nHowever, the common practice has the disadvantage that a functional change in the operating program of the safety controller can be carried out only by the manufacturer of the safety controller himself. If a functional change or an update is desired in the operating program, the user must either send the safety controller to the manufacturer or must request specialist or approved servicing personnel from the manufacturer. This is inconvenient and expensive and, furthermore, may also be disadvantageous in terms of shutdown times for the machine system in which the safety controller is used.\nWhen no safety factors are involved, such as in the case of commercially available personal computers, it is common practice that a user can carry out software updates on his own responsibility by obtaining new software from the manufacturer and loading it onto the personal computer, possibly after instruction. This also applies to so-called operating systems which represent an operating program in the sense of the present invention. According to well established opinion, however, such a procedure is not feasible for safety applications because this would result in the manufacturer of the safety controller losing sole control over the combination of hardware and operating program. In consequence, unchecked combinations of hardware and operating programs would be possible, which would represent a safety risk."}
{"text": "1. Field of the Invention\nThe present invention pertains to the field of advanced sporting equipment design and in particular to a golf club head system for a putter, driver, or iron designed for control of spin resulting from impact between the club head and a golf ball through elastically tailoring normal and tangential impact compliance.\n2. Background Art\nThe present invention pertains to achieving an increase in the accuracy and distance of a golf club (e.g., a driver, putter or iron) through the application of structural design techniques and elastic tailoring of the club and in particular to enhancing or diminishing ball spins. There have been many improvements over the years which have had measurable impact on the accuracy and distance which a golfer can achieve. Typical passive performance improvements such as head shape and volume, weight distribution and resulting components of the inertia tensor, face thickness and thickness profile, face curvatures and CG locations, all pertain to the selection of optimum constant physical and material parameters for the golf club.\nThe impact between the ball and the head can be modeled as an impact between two elastic/deformable bodies each having freedom to translate and rotate in space i.e., full 6 degrees of freedom (DOF) bodies, each having the ability to deform at impact, and each having fully populated mass and inertia tensors. The typical initial condition for this event is a stationary ball and high velocity head impacting the ball at a perhaps eccentric point substantially on or substantially off the face of the club head. The impact results in high forces both normal and tangential to the contact surfaces between the club head and the ball. These forces integrate over time to determine the speed and direction, forming velocity vector and spin vectors of the ball after it leaves the face, hereafter called the impact resultants. These interface forces are determined by many properties including elasticity of the two bodies, material properties and dissipation, surface friction coefficients, body masses and inertia tensors.\nThe present invention pertains to the design of the elastic structural parameters of the head and in particular the attachment between the head body and the face or face insert such that the impact resultants benefit from the elastic/dynamic response of the clubhead under the impact forces. For example the structural design can be such that the face deflections and dynamic response are selected to maximize or minimize ball spin resulting from the impact. There has been much work in the area of elastic tailoring of a golf club head to influence the impact of the head and the ball and the resulting ball flight.\nU.S. Pat. No. 4,498,672 to Bulla issued Feb. 12, 1985 discloses a clubhead designed so that the elastic response of the club in the normal direction is tuned such that it's flexure frequency matches a distortion frequency of the ball. The goal is to increase flight distance by increasing the Coefficient of Restitution (COR).\nU.S. Pat. No. 5,299,807 to Hutin issued Apr. 5, 1994 discloses a clubhead designed with a thin visco-elastic sheet sandwiched between a face and a club head for improving impact performance and feel. There's no mention of spin, but the patent describes an elastically supported face.\nU.S. Pat. No. 5,316,298 to Hutin issued May 31, 1994 discloses a club head designed with a constrained layer visco-elastic damping treatment mounted on the face and or the body for noise tailoring. There's no mention of spin control or control of impact resultants, but the patent discloses an elastically supported face.\nU.S. Pat. No. 5,505,453 to Mack issued Apr. 9, 1996, perhaps the closest to the present invention, discloses several (2) designs for an elastically supported impact plate whose support can be tuned to maximize normal response and exiting ball velocity for a given player. It essentially uses advanced analytical models (1-d) normal impact only to determine the optimal support stiffness in the normal direction to maximize ball velocity after impact. The patent shows two designs each applied to drivers, irons and putters. There's no mention of spin, but the patent discloses an elastically supported face.\nU.S. Pat. No. 5,674,132 to Fisher issued Oct. 7, 1997 discloses a club head designed with an elastically tailored face insert designed to have an desired rebound factor and/or feel/hardness. There's no mention of spin, but the patent discloses an elastically tailored face.\nU.S. Pat. No. 5,697,855 to Aizawar issued Dec. 16, 1997 discloses a clubhead (iron and driver) designed with an elastically supported face insert designed to have a desired damping factor. There's no mention of spin, but the patent discloses an elastically supported face insert.\nU.S. Pat. No. 5,807,190 to Krumme et al. issued Sept. 15, 1998 and U.S. Pat. No. 6,277,033 to Krumme et al. issued Aug. 21, 2001 disclose a clubhead (iron and driver—190, and putter—033) designed with an elastically tailored face comprising a number of pixels each selected for its elastic properties and selectively arranged to give a desired face effect (sweet spot etc). There's no mention of spin, but the patent discloses an elastically tailored face design.\nU.S. Pat. No. 6,001,030 to Delaney et al. issued Dec. 14, 1999 discloses a club head, (putter only) designed with a face insert constructed “with controlled compression”, i.e., a rigid face impact plate elastically supported where the support is designed to provide a certain normal motion behavior depending on impact intensity and/or impact location. There is no mention of spin, but the patent discloses an elastically tailored face design.\nU.S. Pat. No. 6,302,807 to Rohrer issued Oct. 16, 2001 discloses a golf club head (preferably putter) designed with variable energy absorption. It discloses designs for viscoelastic supported faces constructed to maximize dissipation in ideal hits and lower dissipation in off center miss-hits. There's no mention of spin, but the patent discloses an elastically tailored face design.\nU.S. Pat. No. 6,328,661 to Helmstetter et al. issued Dec. 11, 2001 and U.S. Pat. No. 6,478,690 to Helmstetter et al. issued Nov. 12, 2002, “Multiple Material Golf Club Head with a Polymer Insert Base” disclose a golf club head (preferably putter) designed with a polymer face insert of carefully defined hardness and rebound i.e., an elastically tailored insert to effect impact COR and feel.\nU.S. Pat. No. 6,332,849 to Beasley et al. issued Dec. 25, 2001, “Golf Club Driver with Gel Support of Face Wall” discloses a golf club head (preferably driver) designed with a viscoelastic member supporting the face and connected between the center of the face and the back of the hollow body of the clubhead.\nU.S. Pat. No. 6,354,961 to Allen issued Mar. 12, 2002, “Golf Club Face Flexure Control System” discloses a golf club head (preferably driver) designed with a pneumatic piston/cylinder supporting the face and connected between the center of the face and the back of the hollow body of the clubhead. The piston is designed to make contact and change effective stiffness in a predetermined impact velocity range.\nU.S. Pat. No. 6,364,789 to Kosmatka issued Apr. 2, 2002, “Golf Club Head” discloses a golf club head designed with an annular deflection enhancement member disposed between the club head body and a stiff face. The stiffness of the annular member is preferably lower then the face to enhance deflection of the face at impact and increase COR.\nU.S. Pat. No. 6,478,693 to Matsunaga et al. issued Nov. 12, 2002, “Golf Club Head” discloses a golf club head (preferably driver or iron) designed with a variable thickness face with step changes in multiple tiered thickness regions. The centroids of the regions are designed and located to maximize the region of uniformity of strike response—i.e., increase the sweet spot under normal impact.\nU.S. Pat. No. 6,488,594 to Card et al. issued Dec. 3, 2002, “Putter with a consistent Putting Face” discloses a putter designed with a face insert designed to maximize dissipation in ideal hits and lower dissipation in off center miss-hits. There's no mention of spin, but the patent discloses an elastically tailored face design.\nU.S. Pat. No. 6,592,468 to Vincent et al. issued Jul. 15, 2003, “Golf Club Head” discloses a golf club head designed with a viso-elastically supported insert for increasing the damping in vibrations in the club caused by impact.\nU.S. Pat. Nos. 6,595,057 and 6,605,007 to Bissonnette et al. issued Jul. 22, 2003 and Aug. 12 2003, respectively, “Golf Club Head with High Coefficient of Restitution” discloses a golf club with a face whose thickness is tailored to maximize COR. The face has a higher stiffness central zone and a lower stiffness surrounding zone.\nU.S. Pat. No. 6,602,150 to Kosmatka issued Aug. 5, 2003, “Golf Club Striking Plate with Vibration Attenuation” discloses a golf club with a variable thickness face (thicker central portion) on which is disposed a viscoelastic material for face vibration attenuation.\nAll of the aforementioned patents deal with clubhead designs such that the elastic response of the head and face during impact impart a benefit to feel and or COR of the clubhead. None of the aforementioned patents has addressed the design of the elastic/dynamic response of the clubhead so as to effect beneficial control of the ball spin. U.S. Pat. No. 5,193,806 to Burkly issued Mar. 16, 1993, discloses a clubhead designed with a circular shape contact surface to effect spin control, but does not teach the use of clubhead elastic response to achieve this. The face is assumed to be rigid. Numerous patents have attempted to address spin control through surface treatments of the contacting bodies, but none directly address control of spin by elastic/structural design of the clubhead."}
{"text": "1. Technical Field\nMethods and example implementations described herein are directed to an interconnect architecture, and more specifically to systems and methods for implementing visual and/or graphical representation of NoC performance based on outcome of simulation conducted on one or a combination of Network on Chip (NoC) interconnects and/or System on Chip (SoC) architectures.\n2. Related Art\nThe number of components on a chip is rapidly growing due to increasing levels of integration, system complexity and shrinking transistor geometry. Complex System-on-Chips (SoCs) may involve a variety of components e.g., processor cores, DSPs, hardware accelerators, memory and I/O, while Chip Multi-Processors (CMPs) may involve a large number of homogenous processor cores, memory and I/O subsystems. In both SoC and CMP systems, the on-chip interconnect plays a role in providing high-performance communication between the various components. Due to scalability limitations of traditional buses and crossbar based interconnects, Network-on-Chip (NoC) has emerged as a paradigm to interconnect a large number of components on the chip. NoC is a global shared communication infrastructure made up of several routing nodes interconnected with each other using point-to-point physical links denoting connectivity and direction of data flow within the SoC and the NoC.\nMessages are injected by the source and are routed from the source node to the destination over multiple intermediate nodes and physical links. The destination node then ejects the message and provides the message to the destination. For the remainder of this application, the terms ‘components’, ‘blocks’, ‘hosts’ or ‘cores’ will be used interchangeably to refer to the various system components, which are interconnected using a NoC. Terms ‘routers’ and ‘nodes’ will also be used interchangeably. Without loss of generalization, the system with multiple interconnected components will itself be referred to as a ‘multi-core system’.\nThere are several topologies in which the routers can connect to one another to create the system network. Bi-directional rings (as shown in FIG. 1(a)), 2-D (two dimensional) mesh (as shown in FIG. 1(b)) and 2-D Torus (as shown in FIG. 1(c)) are examples of topologies in the related art. Mesh and Torus can also be extended to 2.5-D (two and half dimensional) or 3-D (three dimensional) organizations. FIG. 1(d) shows a 3D mesh NoC, where there are three layers of 3×3 2D mesh NoC shown over each other. The NoC routers have up to two additional ports, one connecting to a router in the higher layer, and another connecting to a router in the lower layer. Router 111 in the middle layer of the example has both ports used, one connecting to the router at the top layer and another connecting to the router at the bottom layer. Routers 110 and 112 are at the bottom and top mesh layers respectively, therefore they have only the upper facing port 113 and the lower facing port 114 respectively connected.\nPackets are message transport units for intercommunication between various components. Routing involves identifying a path composed of a set of routers and physical links of the network over which packets are sent from a source to one or more destination components. Components are connected to one or multiple ports of one or multiple routers; with each such port having a unique ID. Packets carry the destination's router and port ID for use by the intermediate routers to route the packet to the destination components.\nExamples of routing techniques include deterministic routing, which involves choosing the same path from A to B for every packet. This form of routing is independent from the state of the network and does not load balance across path diversities, which might exist in the underlying network. However, such deterministic routing may implemented in hardware, maintains packet ordering and may be rendered free of network level deadlocks. For example, shortest path routing may minimize the latency, as such routing reduces the number of hops from a source to one or more destination(s) and/or reduces the cost of routing a packet from the source to destination(s), wherein the cost of routing depends on bandwidth available between one or more intermediate elements/channels. For this reason, the shortest path may also be the lowest power path for communication between the two components. Dimension-order routing is a form of deterministic shortest path routing in 2-D, 2.5-D, and 3-D mesh networks. In this routing scheme, messages are routed along each coordinates in a particular sequence until the message reaches the final destination. For example in a 3-D mesh network, one may first route along the X dimension until it reaches a router whose X-coordinate is equal to the X-coordinate of the destination router. Next, the message takes a turn and is routed in along Y dimension and finally takes another turn and moves along the Z dimension until the message reaches the final destination router. Dimension ordered routing may be minimal turn and shortest path routing.\nFIG. 2(a) pictorially illustrates an example of XY routing in a two dimensional mesh. More specifically, FIG. 2(a) illustrates XY routing from node ‘34’ to node ‘00’. In the example of FIG. 2(a), each component is connected to only one port of one router. A packet is first routed over the x-axis until the packet reaches node ‘04’ where the x-coordinate of the node is the same as the x-coordinate of the destination node. The packet is next routed over the y-axis until the packet reaches the destination node.\nIn heterogeneous mesh topology in which one or more routers or one or more links are absent, dimension order routing may not be feasible between certain source and destination nodes, and alternative paths may have to be taken. The alternative paths may not be shortest or minimum turn.\nSource routing and routing using tables are other routing options used in NoC. Adaptive routing can dynamically change the path taken between two points on the network based on the state of the network. This form of routing may be complex to analyze and implement.\nA NoC interconnect may contain multiple physical networks. Over each physical network, there may exist multiple virtual networks, wherein different message types are transmitted over different virtual networks. In this case, at each physical link or channel, there are multiple virtual channels; each virtual channel may have dedicated buffers at both end points. In any given clock cycle, only one virtual channel can transmit data on the physical channel.\nThe physical channels are time sliced into a number of independent logical channels called virtual channels (VCs). VCs provide multiple independent paths to route packets, however they are time-multiplexed on the physical channels. A virtual channel holds the state needed to coordinate the handling of the flits of a packet over a channel. At a minimum, this state identifies the output channel of the current node for the next hop of the route and the state of the virtual channel (idle, waiting for resources, or active). The virtual channel may also include pointers to the flits of the packet that are buffered on the current node and the number of flit buffers available on the next node.\nNoC interconnects may employ wormhole routing, wherein, a large message or packet is broken into small pieces known as flits (also referred to as flow control digits). The first flit is the header flit, which holds information about this packet's route and key message level info along with payload data and sets up the routing behavior for all subsequent flits associated with the message. Optionally, one or more body flits follows the head flit, containing the remaining payload of data. The final flit is the tail flit, which in addition to containing the last payload also performs some bookkeeping to close the connection for the message. In wormhole flow control, virtual channels are often implemented.\nThe term “wormhole” plays on the way messages are transmitted over the channels: the output port at the next router can be so short that received data can be translated in the head flit before the full message arrives. This allows the router to quickly set up the route upon arrival of the head flit and then opt out from the rest of the conversation. Since a message is transmitted flit by flit, the message may occupy several flit buffers along its path at different routers, creating a worm-like image.\nBased upon the traffic between various end points, and the routes and physical networks that are used for various messages, different physical channels of the NoC interconnect may experience different levels of load and congestion. The capacity of various physical channels of a NoC interconnect is determined by the width of the channel (number of physical wires) and the clock frequency at which it is operating. Various channels of the NoC may operate at different clock frequencies, and various channels may have different widths based on the bandwidth requirement at the channel. The bandwidth requirement at a channel is determined by the flows that traverse over the channel and their bandwidth values. Flows traversing over various NoC channels are affected by the routes taken by various flows. In a mesh or Torus NoC, there may exist multiple route paths of equal length or number of hops between any pair of source and destination nodes. For example, in FIG. 2(b), in addition to the standard XY route between nodes 34 and 00, there are additional routes available, such as YX route 203 or a multi-turn route 202 that makes more than one turn from source to destination.\nIn a NoC with statically allocated routes for various traffic slows, the load at various channels may be controlled by intelligently selecting the routes for various flows. When a large number of traffic flows and substantial path diversity is present, routes can be chosen such that the load on all NoC channels is balanced nearly uniformly, thus avoiding a single point of bottleneck. Once routed, the NoC channel widths can be determined based on the bandwidth demands of flows on the channels. Unfortunately, channel widths cannot be arbitrarily large due to physical hardware design restrictions, such as timing or wiring congestion. There may be a limit on the maximum channel width, thereby putting a limit on the maximum bandwidth of any single NoC channel.\nAdditionally, wider physical channels may not help in achieving higher bandwidth if messages are short. For example, if a packet is a single flit packet with a 64-bit width, then no matter how wide a channel is, the channel will only be able to carry 64 bits per cycle of data if all packets over the channel are similar. Thus, a channel width is also limited by the message size in the NoC. Due to these limitations on the maximum NoC channel width, a channel may not have enough bandwidth in spite of balancing the routes.\nTo address the above bandwidth concern, multiple parallel physical NoCs may be used. Each NoC may be called a layer, thus creating a multi-layer NoC architecture. Hosts inject a message on a NoC layer; the message is then routed to the destination on the NoC layer, where it is delivered from the NoC layer to the host. Thus, each layer operates more or less independently from each other, and interactions between layers may only occur during the injection and ejection times. FIG. 3(a) illustrates a two layer NoC. Here the two NoC layers are shown adjacent to each other on the left and right, with the hosts connected to the NoC replicated in both left and right diagrams. A host is connected to two routers in this example—a router in the first layer shown as R1, and a router is the second layer shown as R2. In this example, the multi-layer NoC is different from the 3D NoC, i.e. multiple layers are on a single silicon die and are used to meet the high bandwidth demands of the communication between hosts on the same silicon die. Messages do not go from one layer to another. For purposes of clarity, the present disclosure will utilize such a horizontal left and right illustration for multi-layer NoC to differentiate from the 3D NoCs, which are illustrated by drawing the NoCs vertically over each other.\nIn FIG. 3(b), a host connected to a router from each layer, R1 and R2 respectively, is illustrated. Each router is connected to other routers in its layer using directional ports 301, and is connected to the host using injection and ejection ports 302. A bridge-logic 303 may sit between the host and the two NoC layers to determine the NoC layer for an outgoing message and sends the message from host to the NoC layer, and also perform the arbitration and multiplexing between incoming messages from the two NoC layers and delivers them to the host.\nIn a multi-layer NoC, the number of layers needed may depend upon a number of factors such as the aggregate bandwidth requirement of all traffic flows in the system, the routes that are used by various flows, message size distribution, maximum channel width, etc. Once the number of NoC layers in NoC interconnect is determined in a design, different messages and traffic flows may be routed over different NoC layers. Additionally, one may design NoC interconnects such that different layers have different topologies in number of routers, channels and connectivity. The channels in different layers may have different widths based on the flows that traverse over the channel and their bandwidth requirements.\nIn a NoC interconnect, if the traffic profile is not uniform and there is a certain amount of heterogeneity (e.g., certain hosts talking to each other more frequently than the others), the interconnect performance may depend on the NoC topology and where various hosts are placed in the topology with respect to each other and to what routers they are connected to. For example, if two hosts talk to each other frequently and require higher bandwidth than other interconnects, then they should be placed next to each other. This will reduce the latency for this communication, which thereby reduces the global average latency, as well as reduce the number of router nodes and links over which the higher bandwidth of this communication must be provisioned.\nWith the number of on-chip components growing, NoC and SoC being configured to support multiple traffic profiles/transactions/messages having different latency, throughput, and data size characteristics need to be simulated for their performance, and therefore visualization and/or graphical representation of the simulation output including comparison, merging, and conducting other such actions on one or more simulation results becomes necessary in order to evaluate the performance attributes of SoC agents, NoC elements, and/or the NoC channels that form part of the interconnect under varying traffic conditions. There is therefore a need in the art for methods, systems, and non-transitory mediums that can be can be configured for visualization and performance characterization of SoC and/or NoC for one and more transactions."}
{"text": "Endothelial cells line the luminal surface of the vascular bed and are thought to play an active role in the specific proteolytic breakdown of locally deposited fibrin, Todd, J. Pathol. Bacteriol., 78, 281 (1959); Astrup, in Process in Chemical Fibrinolysis and Thrombolysis, Davidson et al. eds., vol. 3, pp. 1-57, Raven Press, New York (1978). The potential of endothelium to initiate and control this process is emphasized by its capacity to synthesize and release plasminogen activators (PAs), Loskutoff et al., Proc. Natl. Acad. Sci. (USA), 74, 3903 (1977); Shepro et al., Thromb. Res., 18, 609 (1980); Moscatelli et al., Cell, 20, 343 (1980); Laug, Thromb. Haemostasis, 45, 219 (1981); Booyse et al., Thromb. Res., 24, 495 (1981), including both tissue-type and urokinase-type molecules, Levin et al., J. Cell Biol., 94, 631 (1982); Loskutoff et al., Blood, 62, 62 (1983). Endothelial cells can also produce inhibitors of fibrinolysis, Loskutoff et al., Proc. Natl. Acad. Sci. (USA), 74, 3903 (1977); Levin et al., Thromb. Res., 15, 869 (1979); Loskutoff et al., J. Biol. Chem. , 256, 4142 (1981); Dosne et al., Thromb. Res., 12, 377 (1978); Emeis et al., Biochem. Biophys. Res. Commun., 110, 392 (1983); Loskutoff et al., Proc. Natl. Acad. Sci. (USA), 80, 2956 (1983); Levin, Proc. Natl. Acad. Sci. (USA), 80, 6804 (1983).\nAlthough these inhibitors probably serve important regulatory roles in controlling the fibrinolytic system of the vascular wall, little is known about their specificity, mode of action, or biochemical nature. The conclusion that these inhibitors are actually synthesized by endothelial cells is obscured somewhat by recent reports that cultured cells can bind and internalize protease inhibitors from serum-containing culture medium, Cohen, J. Clin. Invest., 52, 2793 (1973); Pastan et al., Cell, 12, 609 (1977); Rohrlich et al., J. Cell Physiol., 109, 1 (1981); McPherson et al., J. Biol. Chem., 256, 11330 (1981).\nThe possibility of producing relatively unlimited amounts of tissue-type plasminogen activator (t-PA) by recombinant DNA technology as described in British patent application GB 2,119,804 A, published Nov. 23, 1983, has generated much interest, both clinically and commercially. The conversion of the relatively inactive molecule into an extremely efficient thrombolytic agent by fibrin itself, suggests that t-PA can exist as an active enzyme only when localized to the fibrin-platelet thrombus itself. Thus, t-PA is considered to be a much more specific thrombolytic agent than urokinase-type plasminogen activator and streptokinase.\nThe interactions between t-PA and fibrin have raised the argument that natural inhibitors of t-PA are not necessary to regulate this system; i.e., regulation is achieved through the formation/dissolution of fibrin, and, thus, do not exist. It is clear that the existence of such inhibitors in human blood would complicate attempts to design a specific, efficient, and safe thrombolytic program based upon natural and genetically engineered t-PA. At the very least, calculations such as those of dose, treatment time and efficacy of treatment would be difficult to predict and/or monitor. This problem would be especially acute if inhibitor levels varied from individual to individual.\nThe existence of specific inhibitors of t-PA in plasma is a matter of some dispute, Collen, Thromb. Haemostas., 43, 77 (1980). In fact, it has been reported, Korninger et al., Thromb. Haemostas., 46, 662 (1981), that the activity of t-PA added to plasma had an in vitro half-life of 90 minutes as compared to an in vivo half-life of 2 minutes, Korninger et al., Thromb. Haemostas, 46, 658 (1981). Based upon these observations, those authors concluded that t-PA inhibition by plasma was physiologically unimportant.\nThat conclusion has recently been challenged in Kruithof et al., Prog. in Fibrinolysis, 6, 362 (1983). In Chmielewska et al., Thromb. Res., 31, 427 (1983), direct evidence was recently reported for the existence of a rapid inhibitor of t-PA in plasma. In all cases, this anti-t-PA activity was detected in the plasma of patients with or at risk to develop thrombotic problems; i.e., the very individuals most likely to receive t-PA therapy. This finding may account for the failure of Korninger et al., Thromb. Haemostas., 46, 662 (1981), to detect such an activity since they only examined the plasma of \"normal\" individuals. These reports on t-PA inhibitors represent little more than qualitative descriptions of an \"activity\" detected in the blood of some individuals.\nRecently, an antifibrinolytic agent in cultured bovine endothelial cells was detected, Loskutoff et al., Proc. Natl, Acad. Sci. (USA), 80, 2956 (1983). This inhibitor is a major endothelial cell product and is an inhibitor of plasminogen activator since it san neutralize the activity of both fibrin-independent (urokinase-type) and fibrin-dependent (tissue-type) plasminogen activators (PAs). The observation that human platelets contain an immumologically similar inhibitor, Erickson et al., Haemostasis, 14 (1), 65 (1984) and J. Clin. Invest. 74, 1465 (1984), that is released by them in response to physiologically relevant stimuli, e.g., thrombin, and in parallel with other platelet proteins, e.g., Platelet Factor 4, emphasizes the potential importacute of this inhibitor in human biology. Antiserum to the plasminogen activator inhibitor (PAI) from bovine aortic endothelial cells (BAEs) has been employed to show that the human endothelial PAI as well as that from plasma, serum and platelets are related, i.e., immunologically similar. Erickson et al., Proc. Natl. Acad. Sci. USA, 82, 8710 (1985).\nThe inhibitor found by Loskutoff et al., Proc. Natl. Acad. Sci. (USA), 80, 2956 (1983), was purified from bovine aortic endothelial cell conditioned media by a combination of concanavalin A affinity chromatography and preparative sodium dodecyl sulfate-polyacrylamide gel electrophoresis (SDS-PAGE), and was shown to be a single chain glycoprotein of a molecular weight of 50,000 daltons, having an isoelectric point of 4.5-5 [van Mourik et al., J. Biol. Chem. 259, 14914 (1984)].\nRecent evidence indicates that there are three immunologically distinct plasminogen activator inhibitors (PAIs). The first is that discussed above that is derived primarily from endothelial cells. The second, reported by Astedt et al., Thromb. Haemostasis, 53, 122 (1985) was isolated from placenta. The third, reported by Scott et al., J. Biol. Chem., 260 7029 (1985) is protease nexin.\nThe endothelial cell type PAI differs, in addition to immunologically, from placental PAI and protease nexin in that it inhibits both single chain and two chain tissue-type plasminogen activator (t-PA) as well as urokinase-type plasminogen activator (u-PA), while protease nexin and the placental PAI exhibit substantially no t-PA inhibition at physiological concentrations. Those latter two inhibitors do inhibit u-PA activity at physiological concentrations. Still further, endothelial PAI exhibits beta-mobility when analyzed by agarose zone electrophoresis while the other two PAIs do not. In addition the endothelial cell PAI is stable to low pH values (e.g. pH 3) and SDS (0.1%), while the other two inhibitors are rapidly inactivated by either of these treatments. [van Mourik et al., J. Biol. Chem., 259, 14914 (1984)].\nThe results discussed hereinafter illustrate that the endothelial cell type\nexhibiting , beta-mobility is also present in human placental extracts as is the placenta PAi reported by Astedt et al., Thromb. Haemostasis, 53, 122 (1985). Since two types of PAI are obtainable from placenta, the human\nhereinbefore referred to as of endothelial cell origin will usually be referred to as beta-PAI or endothelial PAI, or endothelial cell type PAI while the PAI first isolated from placenta is referred to as placental-type PAI or placental PAI."}
{"text": "The present invention relates generally to cloth spreading machines and in particular to a method and apparatus for spreading cloth in a uniform manner on a cutting surface.\nCloth spreading machines are widely used in the garment industry and are generally employed to deposit a layer or layers of cloth on a cutting surface. One type of prior art spreader is mounted above a stationary cutting table and is adapted to travel reciprocally above the table, dispensing a layer of cloth from a cloth supply spindle, as it traverses the length of the table. If multiple layers of cloth are to be laid, the spreader will either return to its starting position and dispense another layer of cloth on the table or alternately, a second layer will be deposited during the return travel of the spreader.\nMany of the prior spreaders, include sensors for monitoring the tension in the cloth web. These sensors attempt to control the rate of cloth feed so that the cloth is spread onto the cutting surface uniformly at a predetermined tension, or \"tension free\".\nIn some prior art spreaders, a sensor in the form of a dancer roll is used to monitor cloth tension. The dancer roll is generally pivotally mounted and interposed in the web path so that changes in web tension will cause positional changes in the dancer roll. In one spreader, the dancer roll is operatively connected to a switch that controls the cloth drive; and in another spreader, the dancer roll is connected to a potentiometer which modifies the speed of the cloth drive.\nBecause most dancer rolls are spring biased, a predetermined tension in the web is necessary to overcome the dancer roll force and move it to an equilibrium position. In machines that use dancer rolls, it is difficult, if not impossible, to deposit a uniform layer of cloth on a cutting table, substantially \"tension free\", for the dancer roll itself will produce some tension in the web.\nOther sensors have been suggested which do not contribute to the web tension. These include loop sensors which monitor the droop or sag in the web at a predetermined location. In one suggested spreader, the droop in the web is sensed optically whereas in another spreader, the droop is monitored by a radiant energy sensor. These sensing devices have not been totally satisfactory and in the case of the radiant energy sensor, have been rather expensive.\nA variation in tension along the width of the cloth web is often encountered. Most prior art spreaders do not address this problem, but merely monitor the tension at one position on the web and adjust the cloth feed rate accordingly. Compensation for tension variation across the web is ignored. In the case of patterned fabric, specifically plaids, a variation in tension will skew the orientation of the cloth with respect to the cutting surface and a pattern mismatch will result when the garment is assembled.\nWith the advent of automatic cutting machines, new problems arise and the old problems are compounded. One type of automatic cutting machine is a laser cutter which includes a computer controlled laser head that directs a minute laser beam over a ply of cloth on a cutting surface and cuts out individual garment pieces while simultaneously fusing the edges. The cutting surface is a honey combed conveyor that extends to either side of the laser cutting station. A ply of cloth is laid upon the loading side of the conveyor surface and is then advanced into the cutting station. The cut pieces and the remnants are transported by the conveyor out of the cutting station to an unloading area.\nSpreaders adapted for automatic cutting systems have been proposed. For the laser cutting system described above, a spreader has been proposed which is positioned at the input end of the conveyor and includes a cloth feed roll, an associated drive motor and an optical loop sensor for maintaining a predetermined tension in the web as it is deposited on the conveyor. In the suggested spreader, the cloth feed drive is controlled by the laser cutting machine and dispenses cloth onto the conveyor as it advances toward the cutting station. Because the conveyor speed is quite high, the cloth feed must accelerate quickly to match the speed of the conveyor. Any lag in the cloth feed mechanism will be manifested as an area of excessive tension or stretch in the cloth layer. Garment pieces cut from these stretched areas will not be dimensionally stable. The problem is further aggravated with pattern fabrics, especially plaids, for any nonuniformity is reflected in stripe or pattern skewing. In practice, it has been found that this spreader could not be made to spread cloth consistently on a high speed laser cutting system.\nAnother problem associated with this and other prior suggested spreaders is the inability of the operator adequately to inspect the cloth prior to advancing into the cutting station to insure that it has been spread uniformly. Moreover, even if the operator observes nonuniformity, no provision for rewinding the cloth onto the spreader is provided so that it can be re-spread.\nFinally, it is quite common for remnants to remain on, and by adhered to, the cutting surface after leaving the unloading station. These remnants should be removed because the presence of one will cause distortions in subsequent spreading and cutting operations. Because the prior suggested spreaders are mounted at the input end of the conveyor, it is virtually impossible for the operator to inspect the conveyor prior to cloth spreading, and as a consequence, it is difficult to be certain that all remnants have been removed."}
{"text": "The present invention relates to wireless communication systems in general and, more particularly, to an improved method and apparatus for controlling the power of signals transmitted by base stations and mobile units operating in a CDMA communications system.\nMany wireless communications systems in use today employ a form of spread-spectrum communications technology known as code-division multiple-access, or simply xe2x80x9cCDMAxe2x80x9d. In systems utilizing spread spectrum communications technology such as CDMA, signals transmitted by a base station or mobile unit are spread out over a very wide frequency range using pseudo-random noise sequences, making the signals relatively immune to frequency-dependent interference.\nThe majority of wireless communications systems using CDMA are based on TIA/EIA standard IS-95, which is incorporated by reference herein. IS-95 is known in the industry as a second generation standard. It is anticipated that there will soon be a new industry-wide CDMA standard emerging from various proposals currently under consideration by the International Telecommunications Union (ITU). Such proposals are collectively known as third generation (3G) CDMA proposals and include xe2x80x9ccdma2000 RTT Candidate Submission to ITU-Rxe2x80x9d, which is incorporated by reference herein.\nIn any CDMA system, the mobile units in a given cell act as geographically disparate signal sources which activate at random times. Consequently, it is not possible to synchronize reverse-link transmissions, i.e., transmission from the various mobile units to the cell\"\"s base station. It is therefore impossible for a base station to perform accurate detection (even using pilot-assisted coherent reception as in 3G CDMA) without relying on a feedback mechanism to adjust the transmit power of each mobile unit. As a result, a dynamic method of power control known as closed-loop power control is commonly employed for controlling reverse-link power.\nOn the other hand, forward-link CDMA signals transmitted by a base station and destined for the various mobile units in the cell are designed to be mutually orthogonal and, furthermore, transmission of these signals can be synchronized by the base station. Thus, not only is a transmitted CDMA signal destined for a given mobile unit immune to interference from signals destined for other mobile units, but the mobile unit is able to perform coherent detection with a large processing gain. Accordingly, second generation standards such as IS-95 do not provide for closed-loop power control in the forward-link direction. Nevertheless, it has been found from experience with IS-95 CDMA that significant additional performance improvements can be achieved by using a feedback power control mechanism as in the reverse-link. Accordingly, most 3G CDMA proposals call for the use of closed-loop power control.\nClosed-loop power control consists of a destination (which could be a base station or a mobile unit) measuring the signal-to-interference ratio of a signal received from a source and comparing the measured signal-to-interference ratio with a predetermined target value. If the measured value is greater than the target value, then the power transmitted by the source may be lowered, while if it is less than the target value, then the power transmitted by the source must be increased in order to meet the target. The desired power adjustment is forwarded by the destination back to the source along an appropriate power control subchannel.\nA separate outer-loop power control mechanism is responsible for setting the target signal-to-interference ratio. The guidelines for implementing open-loop power control in the forward-link (3G CDMA) and the reverse-link (IS-95 and 3G CDMA) directions are fairly flexible. Still, thus far, designers have chosen to constrain the target signal-to-interference ratio to a constant value or one which depends on the type of service being offered. While this is satisfactory for situations where the data rate of a service is generally constant and quite low, many emerging applications require high data rates (certain 3G CDMA proposals anticipate that data connections at rates of up to 2 megabits per second per connection should be supported) and are characterized as having a substantially bursty traffic pattern.\nAs a result, if a conventional closed-loop power control algorithm relying on a fixed or service-specific target signal-to-interference ratio is applied in, say, the reverse-link direction, then whenever there is an increase in the instantaneous traffic from a given mobile unit, all the other mobile units will be required to raise their transmit power significantly because of interference from the given mobile unit. The fact that the other mobile units have raised their transmit power causes interference to the given mobile unit, which is again forced to increase its transmit power, and so on. This phenomenon will be repeated at a time scale of the burst duration which makes it a frequently occurring event.\nFrequent increases in transmit power may cause undue stress on the transmitter power amplifier, which may consequently limit the cell capacity. Furthermore, the power control algorithm is required to converge quickly in order to avoid outages during the bit rate transition from a low bit rate to a high bit rate.\nSimilar effects occur in the forward-link direction, although other mobile units within the same cell as the given mobile unit are relatively immune to interference caused by forward-link transmissions involving the given mobile unit. Still, the above described effects can occur as a result of bursty forward-link transmissions occurring in neighbouring cells.\nClearly, the signal-to-interference ratio alone does not appear to be an adequate characterization of the grade of service required of a bursty link. For instance, when the mobile unit is transmitting at a high bit rate, its required performance and required quality of service are likely different from when the bit rate is considerably lower. Thus, there is a need in the industry to provide a method and apparatus for controlling transmitted power by a base station or mobile unit which takes into account the bursty nature of data traffic.\nThe invention can be summarized according to a first broad aspect as a method for use in a communications system wherein the transmit power of a wireless link is adjusted so that link performance meets a target level, the method including dynamically adjusting the target level as a function of the traffic characteristics of the link.\nThe invention can be summarized according to a second broad aspect as a method for use in a closed-loop power control system wherein the transmit power of a source unit communicating with a destination unit across a wireless link is varied in accordance with measured performance and a target performance parameter. The method is one of setting the target performance parameter and is performed at the destination unit. The method includes (a) detecting the start and end of data bursts received from the source unit across the link; and (b) if the performance of the link is adequate, then gradually increasing the target performance parameter when the start of a burst is detected and gradually decreasing the target performance parameter when the end of a burst is detected.\nAccording to another broad aspect, the invention can be summarized as a method of generating power control commands for transmission to a source unit communicating with a destination unit. The method includes (a) measuring an instantaneous performance parameter of signals received from the source unit; (b) determining an instantaneous bit rate, denoted RMINST, of the signals received from the source unit; (c) computing a threshold instantaneous performance parameter as a first function of the measured instantaneous performance parameter, of RMINST, and of at least one target error performance-parameter; and (d) generating a power control command based upon a second function of the measured instantaneous performance parameter and the threshold instantaneous performance parameter.\nPreferably, the measured instantaneous performance parameter is a measured instantaneous signal-to-interference ratio, denoted (S/I)MINST, and the threshold instantaneous performance parameter is a threshold instantaneous signal-to-interference ratio, denoted (S/I)LINST. Preferably, two of the target error performance parameters are a target average signal-to-interference ratio, denoted (S/I)*AVG, and a target maximum outage time, denoted T*MAX.\nPreferably, step (c) includes (c1) integrating (S/I)MINST over a first time window to compute a measured average signal-to-interference ratio, denoted (S/I)MAVG; (c2) integrating RMINST over a second time window to compute an average bit rate, denoted RMAVG; and (c3) among a plurality of most recent values of (S/I)MINST, determining the longest amount of time during which (S/I)MINST was below (S/I)*AVG, said longest amount of time being denoted TMMAX.\nPreferably, the first function is:                               (                      S            I                    )                INST        L            =                                    (                          S              I                        )                    AVG          *                ·        α              ,    where        α    =          {                                                  1              ,                                                                          if                ⁢                                  xe2x80x83                                ⁢                                  T                  MAX                  M                                             greater than                                                 T                  MAX                  *                                ⁢                                  xe2x80x83                                ⁢                or                ⁢                                  xe2x80x83                                ⁢                                                      (                                          S                      I                                        )                                    AVG                  M                                             less than                                                 (                                      S                    I                                    )                                AVG                *                                                                                                                          R                  AVG                  M                                                  R                  INST                  M                                            ,                                                          otherwise              .                                          \nThe invention can be summarized according to another broad aspect as a computer-readable storage medium which, when processed by a computer at a destination unit, executes a sequence of steps to generate a power control command for transmission to a source unit communicating with the destination unit. The steps include (a) computing a threshold instantaneous performance parameter as a first function of a measured instantaneous performance parameter, a measured instantaneous bit rate and at least one target error performance parameter; and (b) generating a power control command based upon a second function of the measured instantaneous performance parameter and the threshold instantaneous performance parameter.\nAccording to still another broad aspect, the invention may be summarized as a statistical power control block for use in a destination unit communicating with at least one source unit. The statistical power control block includes a unit for computing a threshold instantaneous performance parameter as a first function of a measured instantaneous performance parameter, a measured instantaneous bit rate and at least one target error performance parameter; and a unit connected to the computing means, for generating a power control command based upon a second function of the measured instantaneous performance parameter and the threshold instantaneous performance parameter.\nThe invention can also be summarized as a functional unit for generating power control commands. The functional unit includes a maximum selector for receiving a plurality of composite correlation levels and selecting the largest thereamong, as well as an integrator connected to the maximum selector, for receiving the largest composite correlation level and integrating it over a predetermined length of time, thereby to produce a measured instantaneous performance parameter at regular intervals. In addition, the functional unit includes the above-described statistical power control block.\nWhen the invention is implemented, a sudden burst of data results in a lowering of the threshold instantaneous signal-to-interference ratio, until such time that the average signal-to-interference ratio or the maximum outage time are no longer acceptable. This provides a smoothing effect of the interference induced to other users and there may result an increase in cell capacity. Similarly, if the data rate is reduced while the limitations on the average error performance and outage times are being respected, then the average bit rate will be higher than the instantaneous bit rate, which increases the threshold instantaneous signal-to-interference ratio, which further maintains the mobile unit\"\"s transmitted power at a relatively high power level for a longer time interval, followed by a gradual reduction as the average bit rate drops with time."}
{"text": "This invention relates generally to toy water shooters, and more particularly to a novelty toy water shooter that may be concealed under the user\"\"s clothing.\nYoung people of all ages enjoy water fights with toy water guns. As such, it is not surprising there are many, many different types of water shooters available. These run from simple hand-held squirt guns that use trigger-activated pumps to eject water, to more complicated and sophisticated shooters that rely upon pressurized tanks to shoot a stream of water a significant distance.\nA concealed water gun adds an extra dimension of fun to water fights. Among other advantages available when the water gun is hidden, the gun may go unnoticed until it is used to douse its target, and if the shooter is clever enough and the concealment good enough, even after the target has been hit. U.S. Pat. No. 4,997,110 discloses a prior concealed water shooter. While the gun disclosed in the \"\"110 patent may be concealed, it relies upon an electric pump powered by batteries and activated by an electric switch to eject water from the nozzle, and such electrical components add complexity to what is essentially a toy novelty.\nThere is a need therefore for improved toy water shooters, and in particular, concealable water guns.\nA pressurizable bladder is configured for attachment to a user\"\"s arm in a concealed position. The bladder has a refill tube for adding water, and an air inlet that is attached to a hand pump for pressurizing the bladder. A nozzle and valve are fluidly connected to the bladder and a trigger is operable to selectively open and close the valve.\nIn a second illustrated embodiment the pressure bladder is configured for attachment to the user\"\"s waist, and the other components are modified accordingly. The hand pump and bladder-refill tube are connected to the bladder with conduitsxe2x80x94the hand pump conduit may be extended from the bladder down the user\"\"s sleeve and held in one hand, and the nozzle and trigger are held in the user\"\"s other hand."}
{"text": "1. Field of the Invention\nThe present invention relates to heregulin variants, nucleic acid molecules encoding such variants, and related vectors, host cells, pharmaceutical compositions, and methods. In particular, the invention relates to amino acid substitution variants of human heregulin-β1 having an enhanced affinity for the ErbB-3 and ErbB-4 receptors.\n2. Description of the Related Art\nTransduction of signals that regulate cell growth and differentiation is regulated in part by phosphorylation of various cellular proteins. Protein tyrosine kinases are enzymes that catalyze this process. Receptor protein tyrosine kinases are believed to direct cellular growth via ligand-stimulated tyrosine phosphorylation of intracellular proteins. Growth factor receptor protein tyrosine kinases of the class I subfamily include the 170 kilodalton (kDa) epidermal growth factor receptor (EGFR) encoded by the erbB1 gene. erbB1 has been causally implicated in human malignancy. In particular, increased expression of this gene has been observed in more aggressive carcinomas of the breast, bladder, lung, and stomach.\nThe second member of the class I subfamily, p185neu (also called the ErbB-2 receptor or p185HER2), was originally identified as the product of the transforming gene from neuroblastomas of chemically treated rats. The neu (erbB2 or HER2) gene encodes a 185 kDa receptor protein tyrosine kinase.\nAmplification and/or overexpression of the human erbB2 gene correlates with a poor prognosis in breast and ovarian cancers. Slamon et al., Science 235:177-82 (1987); Slamon et al., Science 244:707-12 (1989). Overexpression of erbB2 has been correlated with other carcinomas including carcinomas of the stomach, endometrium, salivary gland, lung, kidney, colon and bladder. Accordingly, in U.S. Pat. No. 4,968,603, Slamon et al. describe and claim various diagnostic assays for determining erbB2 gene amplification or expression in tumor cells. Slamon et al. discovered that the presence of multiple copies of the erbB2 oncogene in tumor cells indicates that the disease is more likely to spread beyond the primary tumor site, and that the disease may therefore require more aggressive treatment than might otherwise be indicated by other diagnostic factors. Slamon et al. conclude that the erbB2 gene amplification test, together with the determination of lymph node status, provides greatly improved prognostic utility.\nA further related gene, called erbB3 (or HER3), which encodes the ErbB-3 receptor (p180HER3) has also been described. See U.S. Pat. No. 5,183,884; Kraus et al., PNAS USA 86:9193-97 (1989); EP Patent Application No. 444,961A1; Kraus et al., PNAS USA 90:2900-04 (1993). Kraus et al. (1989) discovered that markedly elevated levels of erbB3 mRNA were present in certain human mammary tumor cell lines indicating that erbB3, like erbB1 and erbB2, may play a role in human malignancies. Also, Kraus et al. (1993) showed that EGF-dependent activation of the ErbB-3 catalytic domain of a chimeric EGFR/ErbB-3 receptor resulted in a proliferative response in transfected NIH-3T3 cells. Furthermore, these researchers demonstrated that some human mammary tumor cell lines display a significant elevation of steady-state ErbB-3 receptor tyrosine phosphorylation, further implicating this receptor in human malignancies. The role of erbB3 in cancer has been explored by others, and this gene has been found to be overexpressed in breast (Lemoine et al., Br. J. Cancer 66:1116-21 [1992]), gastrointestinal (Poller et al., J. Pathol. 168:275-80 [1992]; Rajkumer et al., J. Pathol. 170:371-78 [1993]; Sanidas et al., Int. J. Cancer 54:935-40 [1993]), and pancreatic cancers (Lemoine et al., J. Pathol. 168:269-73 [1992], and Friess et al., Clinical Cancer Research 1:1413-20 [1995]).\nThe class I subfamily of growth factor receptor protein tyrosine kinases has been further extended to include the ErbB-4 (HER4) receptor, which is the product of the erbB4 (HER4) gene. See EP Patent Application No. 599,274; Plowman et al., PNAS USA 90:1746-50 (1993); and Plowman et al., Nature 366:473-75 (1993). Plowman et al. found that increased erbB4 expression closely correlated with certain carcinomas of epithelial origin, including breast adenocarcinomas. Diagnostic methods for detection of human neoplastic conditions (especially breast cancers) that evaluate erbB4 expression are described in EP Patent Application No. 599,274.\nThe quest for the activator of the erbB2 oncogene has lead to the discovery of a family of heregulin polypeptides. In humans, the heregulin polypeptides characterized thus far are derived from alternate splicing of a single gene which was mapped to the short arm of chromosome 8 by Lee and Wood, Genomics 16:790-91 (1993).\nHolmes et al. isolated and cloned a family of polypeptide activators for the ErbB-2 receptor which they called heregulin-α (HRG-α), heregulin-β1 (HRG-β1), heregulin-β2 (HRG-β2), and heregulin-β3 (HRG-β3). See Holmes et al., Science 256:1205-10 (1992); WO 92/20798; and U.S. Pat. No. 5,367,060. These researchers demonstrated the ability of the purified heregulin polypeptides to activate tyrosine phosphorylation of the ErbB-2 receptor in MCF7 breast tumor cells. Furthermore, the mitogenic activity of the heregulin polypeptides on SK-ER-3 cells (which express high levels of the ErbB-2 receptor) was also demonstrated.\nHeregulins are large multi-domain proteins that are typically expressed as “pro-heregulins.” Pro-heregulins have been shown to undergo proteolytic processing to a mature soluble form (usually of about 44-45 kDa). Processing has been shown to occur intracellularly or at the cell surface. Domains in the soluble form include (in order from the N- to the C-terminus) an immunoglobulin homology (Ig-like) domain, a spacer region rich in glycosylation sites, and a domain similar to a domain found in EGF that is sufficient for ErbB receptor binding and activation. See Barbacci, et al., J. Biol. Chem. 270:9585-89 (1995).\nThe heregulin EGF-like domains are characterized by substantial structural similarities to (Jacobsen et al., Biochemistry 35:3402-17 [1996]), and limited sequence homology with, EGF residues 1-48 (Holmes, et al., supra). Functional similarities between the heregulin EGF-like domains and EGF have been established by data showing that blocks of EGF sequence substituted into heregulin-β1 do not impair binding to cells co-expressing ErbB-3 and ErbB-2. Barbacci et al., supra.\nWhile heregulins are substantially identical in the first 213 amino acid residues, they are classified into two major types, α and β, based on two EGF-like domains that differ in their C-terminal portions. For example, the heregulin-α EGF-like domain differs from that of the β1-isoform by nine substitutions near the C-terminus. The β-isoform has been reported to bind ErbB receptors with approximately eight to 10-fold higher affinity than the α-isoform. Wen et al., Mol. Cell. Biol. 14:1909-19 (1994).\nThe solution structure of the heregulin-α EGF domain has recently been determined at high resolution by NMR: Jacobsen et al., supra; Nagata et al., EMBO J. 10, 3517-3523 (1994). The salient features of this domain include (1) an N-terminal subdomain containing a central three-stranded β-sheet with an intermittent helix and (2) a smaller C-terminal subdomain that contains a short stretch of β-sheet. The EGF domain is stabilized by three disulfide bonds, two in the N-terminal subdomain and one the C-terminal subdomain. The pairing of the six corresponding cysteine residues is conserved in EGF-like domains from all heregulins and from EGF.\nThe 44 kDa neu differentiation factor (NDF), which is the rat equivalent of human HRG, was first described by Peles et al., Cell, 69:205-16 (1992), and Wen et al., Cell, 69:559-72 (1992). Like the human heregulin polypeptides, NDF has an Ig-like domain followed by an EGF-like domain and lacks a N-terminal signal peptide. Subsequently, Wen et al. carried out “exhaustive cloning” to extend the family of NDFs. Wen et al., Mol. Cell. Biol., 14:1909-19 (1994). This work revealed six distinct fibroblastic pro-NDFs. Adopting the nomenclature of Holmes et al., the NDFs were classified as either α or β polypeptides based on the sequences of the EGF-like domains. Isoforms 1 to 4 are characterized on the basis of a variable region between the EGF-like domain and transmembrane domain. Also, isoforms a, b and c are defined based on variable-length cytoplasmic domains. These researchers conclude that different NDF isoforms are generated by alternative splicing and perform distinct tissue-specific functions. See also EP 505 148; WO 93/22424; and WO 94/28133 (discussing NDF).\nFalls et al., Cell 72:801-815 (1993) describe another member of the heregulin family which they call “acetylcholine receptor inducing activity (ARIA) polypeptide.” The chicken-derived ARIA polypeptide stimulates synthesis of muscle acetylcholine receptors. See WO 94/08007. ARIA is a β-type heregulin and lacks the entire spacer region between the Ig-like domain and EGF-like domain of HRG-α and HRGβ1-β3.\nMarchionni et al., Nature 362:312-318 (1993) identified several bovine-derived proteins that they call “glial growth factors (GGFs).” These GGFs share the Ig-like domain and EGF-like domain with the other heregulin proteins described above, but also have an amino-terminal kringle domain. GGFs generally do not have the complete spacer region between the Ig-like domain and EGF-like domain. Only one of the GGFs, GGFII, has an N-terminal signal peptide. See also WO 92/18627; WO 94/00140; WO 94/04560; WO 94/26298; WO 95/32724 (describing GGFs and uses thereof).\nHo et al. describe another member of the heregulin family called “sensory and motor neuron-derived factor (SMDF).” Ho et al., J. Biol. Chem. 270:14523-32 (1995). This protein has an EGF-like domain characteristic of all other heregulin polypeptides but a distinct N-terminal domain. In addition, SMDF lacks both the Ig-like domain and the spacer region found in other heregulin polypeptides. Another feature of SMDF is the presence of two stretches of hydrophobic amino acids near the N-terminus.\nWhile the heregulin polypeptides were first identified based on their ability to activate the ErbB-2 receptor (see Holmes et al., supra), it has been discovered that certain ovarian cells expressing neu (erbB2) and neu-transfected fibroblasts did not bind or crosslink to NDF, nor did they undergo tyrosine phosphorylation in response to NDF. Peles et al., EMBO J. 12:961-71 (1993). This finding indicated that another cellular component was necessary for conferring full heregulin responsiveness.\nCarraway et al. subsequently demonstrated that 125I-rHRG-β1 177-244 bound to NIH-3T3 fibroblasts stably transfected with bovine erbB3 but not to non-transfected parental cells. These researchers also expressed bovine ErbB-3 receptor in insect cells and showed that HRG-β1 177-244 bound to a preparation of ErbB-3 receptor solubilized from these cells. They concluded that ErbB-3 is a receptor for heregulin and mediates phosphorylation of intrinsic tyrosine residues as well as phosphorylation of ErbB-2 receptor in cells that express both receptors. Carraway et al., J. Biol. Chem. 269:14303-05 (1994). Sliwkowski et al. found that cells transfected with erbB3 alone show low affinities for heregulin, whereas cells transfected with both erbB2 and erbB3 show higher affinities. Sliwkowski et al., J. Biol. Chem. 269: 14661-65 (1994).\nPlowman and his colleagues have similarly studied ErbB-4/ErbB-2 receptor activation. They expressed the ErbB2 receptor alone, the ErbB4 receptor alone, or the two receptors together in human T lymphocytes and demonstrated that heregulin is capable of stimulating tyrosine phosphorylation of ErbB-4, but could only stimulate ErbB-2 phosphorylation in cells expressing both receptors. Plowman et al., Nature 336:473-75 (1993).\nThese observations are consistent with the “receptor cross-talking” concept described previously by Kokai et al., Cell 58:287-92 (1989), Stern et al., EMBO J. 7:995-1001 (1988), and King et al., 4:13-18 (1989). These researchers found that binding of EGF to the EGFR resulted in activation of the EGFR kinase domain and cross-phosphorylation of the ErbB-2 receptor. This is believed to be a result of ligand-induced receptor heterodimerization and the concomitant cross-phosphorylation of the receptors within the heterodimer. Wada et al., Cell 61:1339-47 (1990).\nThus, the ErbB receptors are believed to be activated by ligand-induced receptor dimerization. Specifically, heregulins can bind separately to ErbB-3 and ErbB-4 receptors, but not to the ErbB-2 receptor. However, ErbB-2 is required for signaling, and heterodimers containing ErbB-2 in combination with ErbB-3 or ErbB-4 bind heregulins with higher affinity than homodimers containing ErbB-3 or ErbB-4. Plowman et al., Nature 366:473-75 (1993); Sliwkowski et al., J. Biol. Chem. 269:14661-65 (1994).\nThe biological activities of heregulins have been investigated by several groups. For example, Holmes et al. (supra) found that heregulin exerts a mitogenic effect on mammary cell lines (such as SK-SR-3 and MCF-7). Lewis et al. reported that heregulin-β1 stimulated proliferation and enhanced colony formation in soft agar in a number of human breast and ovarian tumor cell lines. Lewis et al., Cancer Research 56:1457-65 (1996). These researchers also showed that ErbB-2 is a critical mediator of heregulin responsiveness.\nPinkas-Kramarski et al. found that NDF (rat heregulin) is expressed in neurons and glial cells in embryonic and adult rat brain and primary cultures of rat brain cells, and suggested that NDF may act as a survival and maturation factor for astrocytes. Pinkas-Kramarski et al., PNAS USA 91:9387-91 (1994). Danilenko et al. reported that the interaction of NDF and the ErbB-2 receptor is important in directing epidermal migration and differentiation during wound repair. Danilenko et al., Abstract 3101, FASEB 8 (4-5):A535 (1994).\nMeyer and Birchmeier analyzed expression of mouse heregulin during embryogenesis and in the perinatal animal using in situ hybridization and RNase protection experiments. Meyer and Birchmeier, PNAS USA 91:1064-68 (1994). These authors conclude, based on expression of this molecule, that heregulin plays a role in vivo as a mesenchymal and neuronal factor. Their findings also indicated that heregulin functions in the development of epithelia.\nFalls et al. (supra) found that chicken ARIA plays a role in myotube differentiation, namely affecting the synthesis and concentration of neurotransmitter receptors in the postsynaptic muscle cells of motor neurons. Corfas and Fischbach demonstrated that ARIA also increases the number of sodium channels in chick muscle. Corfas and Fischbach, J. Neuroscience 13:2118-25 (1993).\nBovine GGFs have been reported to be mitogenic for Schwann cells. See, e.g., Brockes et al., J. Biol. Chem. 255:8374-77 (1980); Lemke and Brockes, J. Neurosci. 4:75-83 (1984); Brockes et al., J. Neuroscience 4:75-83 (1984); Brockes et al., Ann. Neurol. 20:317-22 (1986); Brockes, Methods in Enzym. 147:217-225 (1987); Marchionni et al., supra. Schwann cells provide myelin sheathing around the axons of myelinated neurons and thus play an important role in the development, function and regeneration of peripheral nerves. The implications of this role from a therapeutic standpoint have been addressed by Levi et al., J. Neuroscience 14:1309-19 (1994). Levi et al. discussed the potential for construction of a cellular prosthesis including Schwann cells that could be transplanted into areas of damaged spinal cord. Methods for culturing Schwann cells ex vivo have been described. See WO 94/00140; Li et al., J. Neuroscience 16:2012-19 (1996).\nGGFII has been shown to be mitogenic for subconfluent quiescent human myoblasts, and differentiation of clonal human myoblasts in the continuous presence of GGFII results in greater numbers of myotubes after six days of differentiation. Sklar et al., J. Cell Biochem., Abst. W462, 18D, 540 (1994); see also WO 94/26298.\nThe relationship between the structure and function of new proteins can be investigated using any of a variety of available mutational analysis techniques. Examples of such techniques include alanine scanning mutagenesis and phagemid display. Alanine scanning can be used to identify active residues (i.e., residues that have a significant effect on protein function) in a protein or protein domain. For example, Cunningham and Wells used alanine scanning to identify residues in human growth hormone that were important for binding its receptor. Cunningham and Wells, Science 244:1081-85 (1999). In alanine scanning, a gene encoding the protein or domain to be canned is inserted into an expression vector, and mutagenesis is carried out to generate a series of vectors that encode proteins or domains in which sequential residues are converted to alanine. The encoded proteins or domain are expressed from these vectors, and the activities of the alanine-substituted variants are then tested to identify those with altered activity. An alteration in activity indicates that the residue at the alanine-substituted position is an active residue.\nPhagemid display was developed to allow the screening of a large number of variant polypeptides for a particular binding activity. Smith and Parmley demonstrated that foreign peptides can be “displayed” efficiently on the surface of filamentous phage by inserting short gene fragments into gene III of the fd phage. Smith, Science 228:1315-17 (1985); Parmley and Smith, Gene 73:305-18 (1985). The gene III coat protein is present in about five copies at one end of the phage particle. The modified phage were termed “fusion phage” because they displayed the foreign peptides fused to the gene III coat protein. As each fusion phage particle displayed approximately five copies of the fusion protein, this mode of phage display was termed “polyvalent display.”\nScott et al. and Cwirla et al. showed that fusion phage libraries could be screened by sequential affinity selections known as “panning.” Scott et al., Science 249:386-90 (1990); Cwirla et al., PNAS USA 87:6378-82 (1990). However, early efforts to select high affinity fusion phage failed, presumably due to the polyvalence of the phage particles. This problem was solved with the development of a “monovalent” phage display system in which the fusion protein is expressed at a low level from a phagemid and a helper phage provides a large excess of wild-type coat protein. Bass et al., Proteins 8:309-14 (1990); Lowman et al., Biochem. 30:10832-38 (1991). Monovalent phage display can be used to generate and screen a large number of variant polypeptides to isolate those that bind with high affinity to a target of interest."}
{"text": "The present invention relates generally to the field of electrical control and monitoring systems, and more particularly to a system and method that integrates functions control and monitoring into a human machine interface.\nA wide range of systems are known and are currently in use for controlling and monitoring processes, particularly in the industrial context. Such processes may generally include a large number of components, such as pumps, valves, conveyors, material handling and machining systems, and so forth. In most applications a significant number of the components are operated by prime movers such as electric motors. These may be manually operated and inspected, but are often more effectively controlled through programmed systems including protective components as well as activating and control components.\nOne type of industrial control center that has been developed over recent years in generally referred to as a motor control center. Such systems may provide for highly integrated control of a large number of devices and can be equipped for remote control and reporting functions. Moreover, such remote control, typically via industrial control networks, is often highly desirable because it permits locating the protective and power control components roughly in the vicinity of the controlled equipment, which may be quite distant from the more centralized control or monitoring facility or room where human operators are based.\nChallenges faced in implementing and maintaining complex control systems include the configuration of the control system and the representation of the system in a manner that can be quickly and easily mastered by human operators. Moreover, during operation, the representations of the system provided to the human operators, and the meaning of the various controls at the operators' disposal are of considerable importance insomuch as they permit the operator to make efficient and timely decisions based upon accurate understanding of the current and future conditions of the system. Planning, programming and configuration of such systems and networks is, however, quite difficult.\nIn addition to the programming concerns, operating the modem industrial control and monitoring systems often involves different programs that are utilized by the system. These separate programs each require extensive programming to provide information that allows the operator to monitor and control various processes. For instance, one program may include detailed information relating to individual programmable devices, such as manuals and drawings, while another program may provide functionality that enables an operator to view more the complete system that is operating the processes.\nThere is a need, therefore, for an improved technique which may be employed to more effectively present and access information relating to an overall process or system in a human machine interface, while providing detailed and current information on specific components of interest upon demand by an operator. There is a particular need for a technique which will provide greater uniformity and consistency for the operation of the processes by integrating the information of various levels within the modern industrial control and monitoring systems, reducing development and programming time associated with the various programs, and enabling the operator to perform a greater level of diagnostics from a single interface."}
{"text": "Manufactured products often contain orifices and cavities or other hollow parts that result from the manufacturing process and/or that are designed into the product for various purposes, such as weight reduction. Automotive vehicles, for example, include several such orifices and cavities throughout the vehicle, including in the vehicle's structural pillars and in the sheet metal of the vehicle doors. It is often desirable to seal such orifices and cavities so as to minimise noise, vibrations, fumes, dirt, water and the like from passing from one area to another within the vehicle by means of sealing members or baffle elements built into the orifice or cavity. Likewise, such members or elements often fulfill an additional task of reinforcing the hollow structure of the manufactured product, e.g. automotive part, so much that it becomes more resistant to mechanical stress but still maintains the low weight advantage of the hollow structure.\nSuch elements used for sealing, baffling or reinforcing often consist of a carrier, made of plastic, metal, or another rigid material, and one or more layers of a thermoplastic material attached to it which is able to expand its volume when heat or another physical or chemical form of energy is applied. With such a design, it is possible to insert the baffle or reinforcement element into the hollow part of the structure during the manufacture process but also to leave the inner walls of the structure still accessible (or the cavities passable) by e.g. a liquid. For example, during the manufacture process of a vehicle, the hollow parts of a metal frame can still be largely covered by an electrocoating liquid while the baffle or reinforcement elements are already inserted, and afterwards during a heat treatment step, the expandable thermoplastic material layers of the baffle or reinforcement expand to close the cavities as intended.\nThe development of such baffles or reinforcement elements has led to highly advanced systems, where the expandable material is able to increase its volume by up to 2000% or more, forming a foam-like structure filling the cavities and adhering to the walls of the structure intended to be sealed, baffled, or reinforced. Especially in automotive manufacturing, this has led to considerable weight reduction and excellent dampening of noise or vibrations in the car body.\nWith advanced materials that are able to expand even more, it is possible to reduce initial material mass and therefore contribute further to weight reduction and cost efficiency.\nHowever, increasingly small portions of expandable material can also cause problems in manufacturing of the baffle or reinforcement elements. As such elements (or at least the expandable material layers) are normally produced by injection molding or extrusion, the feeding of very small sections gives rise to higher processing demands, such as higher injection pressures, or causes quality issues, such as flashing, short shots (incompletely molded parts) or material degradation.\nIt is thus desirable to obtain a way of manufacturing baffle or reinforcement elements with less initial mass of expandable material, but without the problems connected to very small initial volumes of that material."}
{"text": "1. Field of the Invention\nThe present invention generally relates to executing instructions in a processor. Specifically, this application is related to increasing the efficiency of a processor executing branch instructions.\n2. Description of the Related Art\nModern computer systems typically contain several integrated circuits (ICs), including a processor which may be used to process information in the computer system. The data processed by a processor may include computer instructions which are executed by the processor as well as data which is manipulated by the processor using the computer instructions. The computer instructions and data are typically stored in a main memory in the computer system.\nProcessors typically process instructions by executing the instruction in a series of small steps. In some cases, to increase the number of instructions being processed by the processor (and therefore increase the speed of the processor), the processor may be pipelined. Pipelining refers to providing separate stages in a processor where each stage performs one or more of the small steps necessary to execute an instruction. In some cases, the pipeline (in addition to other circuitry) may be placed in a portion of the processor referred to as the processor core. Some processors may have multiple processor cores, and in some cases, each processor core may have multiple pipelines. Where a processor core has multiple pipelines, groups of instructions (referred to as issue groups) may be issued to the multiple pipelines in parallel and executed by each of the pipelines in parallel.\nAs an example of executing instructions in a pipeline, when a first instruction is received, a first pipeline stage may process a small part of the instruction. When the first pipeline stage has finished processing the small part of the instruction, a second pipeline stage may begin processing another small part of the first instruction while the first pipeline stage receives and begins processing a small part of a second instruction. Thus, the processor may process two or more instructions at the same time (in parallel).\nProcessors typically provide conditional branch instructions which allow a computer program to branch from one instruction to a target instruction (thereby skipping intermediate instructions, if any) if a condition is satisfied. If the condition is not satisfied, the next instruction after the branch instruction may be executed without branching to the target instruction. Typically, the outcome of the condition being tested is not known until the conditional branch instruction is executed and the condition is tested. Thus, the next instruction to be executed after the conditional branch instruction may not be known until the branch condition is tested.\nWhere a pipeline is utilized to execute instructions, the outcome of the conditional branch instruction may not be known until the conditional branch instruction has passed through several stages of the pipeline. Thus, the next instruction to be executed after the conditional branch instruction may not be known until the conditional branch instruction has passed through the stages necessary to determine the outcome of the branch condition. In some cases, execution of instructions in the pipeline may be stalled (e.g., the stages of the pipeline preceding the branch instruction may not be used to execute instructions) until the branch condition is tested and the next instruction to be executed is known. However, where the pipeline is stalled, the pipeline is not being used to execute as many instructions in parallel (because some stages before the conditional branch are not executing instructions), causing the benefit of the pipeline to be reduced and decreasing overall processor efficiency.\nIn some cases, to improve processor efficiency, branch prediction may be used to predict the outcome of conditional branch instructions. For example, when a conditional branch instruction is encountered, the processor may predict which instruction will be executed after the outcome of the branch condition is known. Then, instead of stalling the pipeline when the conditional branch instruction is issued, the processor may continue issuing instructions beginning with the predicted next instruction.\nHowever, in some cases, the branch prediction may be incorrect (e.g., the processor may predict one outcome of the conditional branch instruction, but when the conditional branch instruction is executed, the opposite outcome may result). Where the outcome of the conditional branch instruction is mispredicted, the predicted instructions issued subsequently to the pipeline after the conditional branch instruction may be removed from the pipeline and the effects of the instructions may be undone (referred to as flushing the pipeline). Then, after the pipeline is flushed, the correct next instruction for the conditional branch instruction may be issued to the pipeline and execution of the instructions may continue. Where the outcome of a conditional branch instruction is incorrectly predicted and the incorrectly predicted group of instructions is flushed from the pipeline, thereby undoing previous work done by the pipeline, the efficiency of the processor may suffer.\nAccordingly, what is needed is an improved method and apparatus for executing conditional branch instructions and performing branch prediction."}
{"text": "The present invention relates in general to a bag packing apparatus, and more particularly, an apparatus for automatically opening, packing and sealing, successively, conveyed paper or plastic bags or the like.\nIn a heretofore known bag packing apparatus as shown in FIG. 1, wherein three-sided sealed bags 1, made of paper or plastic films with mouths 2 along their top edge, are charged with contents by means of a charging device 3 and tightly sealed 4, it was a common practice to convey the bags 1 successively in the direction of their width by means of chain conveyors 5. However, such apparatuses in the prior art have the following shortcomings, and consequently, many problems remain unsolved for establishing common availability for different bag sizes and for speeding up processing capability. For example, when the sidth of the bags changed it was necessary to change the mounting distance in the widthwise direction between the bag gripping claws 6 within the conveyor pitch P.sub.1 for each attachment. This resulted in a very complex structure and adjustment requirements for the machine; size change was not easy. In addition, in order to enhance the processing capability, the conveying speed had to be raised to a higher speed due to the fact that the bags were conveyed in the widthwise direction, and especially when the contents of the bags were liquid, there was a disadvantage in that the liquid would squeeze out of the top edge of the bag due to acceleration, resulting in further adverse effects upon sealing."}
{"text": "Electroacoustic transducers are capable of converting electric energy to acoustic energy and vice versa. Electroacoustic receivers typically convert electric energy to acoustic energy through a motor assembly having a movable armature. Typically, the armature has one end that is free to move while the other end is fixed to a housing of the receiver. The assembly also includes a drive coil and one or more magnets, both capable of magnetically interacting with the armature. The armature is typically connected to a diaphragm near its movable end. When the drive coil is excited by an electrical signal, it magnetizes the armature. Interaction of the magnetized armature and the magnetic fields of the magnets causes the movable end of the armature to vibrate. Movement of the diaphragm connected to the armature produces sound for output to the human ear. Examples of such transducers are disclosed in U.S. Pat. Nos. 3,588,383, 4,272,654 and 5,193,116.\nThe sound pressure output of a receiver is created by the travel, or deflection, of the armature when it vibrates. Maximum deflection of the moving armature creates maximum sound pressure output for a given armature geometry. The maximum deflection of an armature is limited by the magnetic saturation of the armature, which is governed by the maximum magnetic flux that the armature geometry can allow to pass therethrough. Therefore, the magnetic flux must be increased in order to increase the sound pressure output. The magnetic flux is limited by material type and cross-sectional area of the armature. Although an increase in the cross-sectional area causes a proportional increase in the maximum magnetic flux, the relative stiffness of the armature increases as well. Thus, merely increasing the cross-sectional area of the armature geometry does not provide a significant improvement in the maximum deflection of the armature.\nThe present invention addresses these and other problems."}
{"text": "Nearly one in seven people in the United States suffer from some type of chronic sleep disorder, and only 50% of people are estimated to get the recommended seven to eight hours of sleep each night. It is further estimated that sleep deprivation and its associated medical and social costs (loss of productivity, industrial accidents, etc.) exceed $150 billion per year. Excessive sleepiness can deteriorate the quality of life and is a major cause of morbidity and mortality due to its role in industrial and transportation accidents. Sleepiness further has undesirable effects on motor vehicle driving, employment, higher earning and job promotion opportunities, education, recreation, and personal life.\nPrimary sleep disorders affect approximately 50 million Americans of all ages and include narcolepsy, restless legs/periodic leg movement, insomnia, and most commonly, sleep apnea. Sleep apnea is defined as the cessation of breathing during sleep. The three major types of sleep apnea are obstructive sleep apnea (OSA), central sleep apnea (CSA), and complex sleep apnea (CompSA). Of these three, CSA is rare, while OSA is the most common. CompSA is a relatively new disease state that manifests itself after therapy is applied. Patients with CompSA are characterized by the emergence of new CSA events after the application of Continuous Positive Airway Pressure (CPAP). OSA's prevalence in society is comparable with diabetes, asthma, and the lifetime risk of colon cancer. OSA is grossly under diagnosed; an estimated 80-90% of persons afflicted have not received a clinical diagnosis. OSA is characterized by repetitive pauses in breathing during sleep due to the obstruction and/or collapse of the upper airway (throat), usually accompanied by a reduction in blood oxygen saturation, and often followed by an awakening to breathe (an apnea event). Respiratory effort continues during the episodes of OSA. Multiple episodes of apnea may occur in one night, causing sleep disruption. CSA is a neurological condition causing cessation of all respiratory effort during sleep, usually with corresponding decreases in blood oxygen saturation. In contrast to OSA, where there is respiratory effort from the brain stem but a physical blockage prevents inhalation of oxygen, in CSA the brainstem center controlling breathing shuts down, resulting in no respiratory effort and no breathing. The subject is aroused from sleep by an automatic breathing reflex. Frequent activation of the reflex results in very little sleep for the subject. The neurological mechanism behind CSA is very different from the physical cause of OSA. Although the effects of CSA and OSA are highly similar, effective treatment can differ. CompSA can be thought of as a combination of OSA and CSA. As mentioned before, CompSA is characterized by an emergence of CSA events after CPAP initiation.\nMedications, hygiene, or some physical form of therapy can be used to treat apneas. Treatment of OSA and CSA vary substantially, which makes a proper diagnosis of the correct type of sleep apnea (OSA, CSA, or CompSA) critical for an effective treatment. Apnea treatment is provided based on the type of apnea, and can be adjusted by re-testing the subject at some later time to determine whether the condition or the symptoms have been alleviated. The most common method of treating OSA is continuous positive airway pressure (CPAP) and positive airway pressure (PAP) devices applied to the subject's airway to force the subject to breathe. When using a simple CPAP device to treat OSA, the air pressure acts as a splint, holding the airway open and reducing or removing the obstruction. The optimal pressure is determined by a sleep technician during a single titration night. The sleep technician manually adjusts the device to deliver a minimum pressure sufficient to force the airway open and reduce the number of apneas. Once the optimal pressure is determined, the device is programmed to consistently provide this pressure, and the patient is sent home.\nSlightly more advanced PAP devices automatically adjust the air pressure based on sensors built into the device. The sensors measure gas flow, pressure, or other fluid characteristics in the device or its mask, and adjust the delivered pressure based on various algorithms known in the art. These auto-PAP devices rely on the single physiological variable (the measured or estimated fluid characteristics) to predict or detect an apnea event.\nNone of the devices on the market can be used to adjust the gas flow delivered to a subject based on a comprehensive evaluation of the subject's current physiological state or the subject's current symptoms. Further, none of the current devices can use a rich data set to predict or detect apnea and provide appropriate treatment. Still further, none of the current CPAP or PAP titration methods use a rich set of data over single or multiple nights to set the optimal pressure and other parameters. Still even further, none of the current devices can be used to automatically adjust a treatment device based on the subject's physiological signals. Still even further, none of the current titration devices can be used in the subject's home.\nIt is therefore an object of the present invention to adjust the treatment gas flow or pressure delivered to the subject based on the subject's current physiological state or symptoms. It is another object of the present invention to use a rich data set over multiple nights to titrate the CPAP treatment. It is another object of the present invention to use a closed-loop or partially closed-loop system to automatically titrate the CPAP treatment based on the subject's physiological signals. It is still another object of the present invention to treat a subject's apnea in a predictive manner. It is still another object of the present invention to provide a system or method of treating a subject's apnea using the subject's physiological signals. It is still another object of the present invention to provide a device and method of titration in the subject's home. It is still another object of the present invention to provide a device and method of titration in the hospital's acute or sub-acute settings, such as for postoperative management of care."}
{"text": "1. Field of the Invention\nThis invention relates to electronic circuits and, more particularly, to clock generation and timing circuits.\n2. Description of the Related Art\nThe following descriptions and examples are given as background only.\nPhase locked loops (PLLs) are commonly used for data and telecommunications, frequency synthesis, clock recovery, and similar applications. In some cases, a PLL may be used in the I/O interfaces of digital integrated circuits to hide clock distribution delays and to improve overall system timing. In other cases, a PLL may be used as a clock multiplier for downstream circuit components. In one example application, an input clock of 10 Mhz can be multiplied by the PLL to yield a 1000 Mhz output signal, which may then be used for clocking internal circuit components. Ideally, this input (or source) clock multiplication could result in an output that is in perfect phase alignment with the input clock. However, in practice this is often not the case.\nA typical PLL device 100 is shown in FIG. 1 as including an optional reference divider (div M) 120, a phase frequency detector (PFD) 130, a charge pump 140, a low pass filter 150, a voltage controlled oscillator (VCO) 160, and an optional frequency divider (div N) 170.\nDuring operation, PLL circuit 100 receives a reference clock signal (FREF) from an external source (e.g., a crystal oscillator) 110. The phase frequency detector compares the reference signal (FREF) to a feedback signal (FVCO) generated by components within the PLL circuit. More specifically, PFD 130 detects differences in frequency and/or phase between the reference and feedback clock signals, and generates compensating “up” and “down” signals in response thereto. The particular control signals generated depend on whether the feedback clock signal is lagging or leading the reference clock signal in frequency or phase. The up/down control signals are passed through charge pump 140 and filter 150 to integrate the control signals into a control voltage, which is sent to the VCO. The voltage-controlled oscillator converts the voltage information into one or more output frequencies (FVCO). One of these output frequencies may then be sent back to the PFD via a feedback loop.\nIn some cases, frequency divider 170 and reference divider 120 may be included for adjusting the frequencies of the feedback and reference clock signals, respectively. For example, frequency divider 170 may be used for dividing the frequency of a VCO output signal (FOUT) to produce a divided down feedback signal (FOUT/N), while reference divider 120 is used for dividing the frequency of the external clocking signal to produce a divided down reference signal (FREF/M), which is similar or dissimilar to the divided down feedback signal. In such cases, PLL 100 may function as a clock or frequency multiplier. However, dividers 120 and/or 170 may not be included in all cases.\nIn conventional PLL devices, high-performance inductor-capacitor (LC) VCOs are often used to create low-noise, high-speed PLLs. For example, LC-type oscillators are often used in PLL designs tailored for wireless and low power applications, as well as other applications requiring precise timing. Unfortunately, LC-type oscillators have a tight frequency range, which is sensitive to variations in process, voltage and temperature (PVT). In some cases, variations in PVT may cause the VCO frequency range to shift outside of a target range, thereby limiting the usable frequency range of the PLL device. For this reason, various solutions have been proposed for extending the usable frequency range of a PLL device employing an LC-type oscillator.\nOne solution to the above-mentioned problem is to create or use a VCO with a wide frequency range. For example, a wider frequency range can be obtained by increasing the gain of an LC-type oscillator, or by using a completely different oscillator (e.g., a ring oscillator) with an inherently wider frequency range. Unfortunately, large VCO gains are undesirable because of their sensitivity to noise. In addition, although a wide frequency range may be easily obtained when the ring oscillator is employed as a VCO, the ring oscillator is not without limitations and usually demonstrates poorer phase-noise performance than the LC-type oscillator.\nAnother solution is to use multiple LC-type oscillators within the PLL device. For example, a first VCO may be used for generating frequencies within a 2-2.45 GHz range, while a second VCO is used for generating frequencies within a 2.45-2.5 Ghz range. Additional or alternative VCOs may be used for generating frequencies within other ranges. Unfortunately, the second solution may be undesirable in many applications, due to the relatively large die area consumed by the additional VCO(s).\nIn yet another solution, an LC-type VCO may be calibrated during the manufacturing test process to shift the VCO operating frequency into a desired range. For example, test circuitry may be used for measuring the max VCO frequency, measuring the min VCO frequency and calculating an average of the two. The test circuitry may then be used for adjusting the programmable bits supplied the LC-type VCO until the measured value(s) equal a set of target value(s). Unfortunately, the third conventional solution is too slow (i.e., adds additional test time) and does not compensate for the environmental conditions that the device will actually be used in (i.e., the method does not account for power supply and temperature variations, only for process variations).\nIn yet another solution, an LC-type VCO may be calibrated during operation of the PLL device to shift the VCO operating frequency into a desired range. One such solution is described in a paper entitled “A Delta-Sigma Fractional-N Synthesizer using a Wide-Band Integrated VCO and a Fast AFC Technique for GSM/GPRS/WCDMA Applications,” and published in the July 2004 issue of IEEE Journal of Solid-State Circuits (JSSCC), vol. 39, no. 7, pgs. 1164-1169. In this solution, a pair of switches are used for disconnecting the low pass filter from the VCO (i.e., opening the loop) to enter a VCO calibration mode. To achieve calibration, digital frequency counters are used for counting the number of reference and feedback clock pulses supplied thereto. A comparator is then used to determine the proper VCO frequency by comparing the number of feedback pulses to the number of reference pulses, while a state machine is used for adjusting the programmable bits supplied to the VCO. Unfortunately, the frequency counters and state machine used in the fourth solution are too slow for many high-speed applications (e.g., it may take about 650 clock cycles to complete the calibration using the method describe above). In addition, the switches used for disconnecting the low pass filter from the VCO introduce a series resistance into the PLL signal path. This is undesirable because it tends to alter the behavior of the filter.\nTherefore, a need remains for an improved calibration solution that may be used for extending the usable frequency range of a PLL device. Preferably, the improved solution would allow frequency calibration and compensation of a narrow band PLL over a wide range of actual environmental conditions, including process, voltage, and temperature. Even more preferably, the improved solution would provide a circuit and method for calibrating an LC-type oscillator during operation of a PLL device, wherein said calibration is performed without opening the loop and with much greater speed than possible with conventional solutions."}
{"text": "I. Field\nThe present invention relates generally to electronics circuits, and more specifically to current sources and active circuits.\nII. Background\nCurrent sources are widely used to provide current for various circuits such as amplifiers, buffers, oscillators, and so on. Current sources may be used as bias circuits to provide bias currents, active loads to provide output currents, and so on. Current sources are often fabricated on integrated circuits (ICs) but may also be implemented with discrete circuit components.\nAs IC fabrication technology continues to improve, the size of transistors continues to shrink. The smaller transistor size enables more transistors and thus more complicated circuits to be fabricated on an IC die or, alternatively, a smaller die to be used for a given circuit. The smaller transistor size also supports faster operating speed and provides other benefits.\nComplementary metal oxide semiconductor (CMOS) technology is widely used for digital circuits and many analog circuits. A major issue with shrinking transistor size in CMOS is leakage current, which is the current passing through a transistor when it is turned off. A smaller transistor geometry results in higher electric field (E-field), which stresses a transistor and causes oxide breakdown. To decrease the E-field, a lower power supply voltage is often used for smaller geometry transistors. However, the lower supply voltage also increases the propagation delay of the transistors, which is undesirable for high-speed circuits. To reduce the delay and improve operating speed, the threshold voltage (Vt) of the transistors is reduced. The threshold voltage determines the voltage at which the transistors turn on. However, the lower threshold voltage and smaller transistor geometry result in higher leakage current.\nLeakage current is more problematic as CMOS technology scales smaller. This is because leakage current increases at a high rate with respect to the decrease in transistor size. Leakage current can impact the performance of certain circuits such as phase lock loops (PLLs), oscillators, digital-to-analog converters (DACs), and so on.\nSome common techniques for combating leakage current include using high threshold voltage (high-Vt) transistors and/or larger transistor sizes (e.g., longer gate lengths). High-Vt transistors may impact circuit performance (e.g., slower speed) and typically require an additional mask step in the IC fabrication process. Larger-size transistors are marginally effective at combating leakage current since (1) leakage current is a relatively weak function of channel length and (2) there are practical limits on how long the channel length may be extended. Both of these solutions may thus be inadequate for certain circuits.\nThere is therefore a need in the art for a current source with low leakage current and good performance."}
{"text": "The semiconductor industry has continually sought ways to produce memory devices with an increased number of memory cells per memory die. In non-volatile memory (e.g., NAND flash memory), one way to increase memory density is by using a vertical memory array, which is also referred to as a three-dimensional (3-D) memory array. One type of vertical memory array includes semiconductor pillars that extend through openings (e.g., holes) in layers of conductive material (also referred to as word line plates or control gate plates), with dielectric materials at each junction of the semiconductor pillars and the conductive materials. Thus, multiple transistors can be formed along each pillar. Vertical memory array structures enable a greater number of transistors to be located in a unit of die area by building the array upwards (e.g., vertically) on a die, as compared to structures with traditional planar (e.g., two-dimensional) arrangements of transistors.\nVertical memory arrays and methods of forming them are described in, for example: U.S. Patent Application Publication No. 2007/0252201 of Kito et al., now U.S. Pat. No. 7,936,004, issued May 3, 2011; Tanaka et al., “Bit Cost Scalable Technology with Punch and Plug Process for Ultra High Density Flash Memory,” Symposium on VLSI Technology Digest of Technical Papers, pp. 14-15 (2007); Fukuzumi et al., “Optimal Integration and Characteristics of Vertical Array Devices for Ultra-High Density, Bit-Cost Scalable Flash Memory,” IEDM Technical Digest, pp. 449-52 (2007); and Endoh et al., “Novel Ultrahigh-Density Flash Memory with a Stacked-Surrounding Gate Transistor (S-SGT) Structured Cell,” IEEE Transactions on Electron Devices, vol. 50, no. 4, pp. 945-951 (April, 2003).\nConventional vertical memory arrays require an electrical connection between the conductive materials (e.g., word line plates or control gates) and access lines (e.g., word lines) so that memory cells in the 3-D array may be uniquely selected for writing or reading functions. One method of forming an electrical connection includes forming a so-called “stair-step” structure at the edge of the conductive materials. FIGS. 1A through 1D show one conventional method of creating a stair-step structure 10 in a stack of conductive materials 12. As shown in FIG. 1A, conductive materials 12 are separated by insulating materials 14 between the conductive materials 12. A mask 16 (e.g., photoresist material) is formed over the topmost insulating material 14 and patterned to expose a portion of the insulating material 14a, the exposed portion having a width of one so-called “step” of the stair-step structure 10 to be formed. An anisotropic etch 18, such as a reactive ion etch (RIE) or other dry etch, is performed to remove the insulating material 14a at the portion exposed through the mask 16. The pattern in the insulating material 14a is then transferred to the conductive material 12a. The exposed insulating material 14a is removed by one dry etch process that stops on the conductive material 12a, and the exposed conductive material 12a is then removed by another dry etch process that stops on the insulating material 14b. Next, the mask 16 is reduced in size by removing a portion of the mask (also known as “trimming”), such as by isotropic etching, to expose another portion of the insulating material 14a, as shown in FIG. 1B.\nThe process is repeated by subjecting the structure to an anisotropic etch 18, including removing exposed portions of the two insulating materials 14a and 14b and subsequently removing exposed portions of the two conductive materials 12a and 12b. As shown in FIG. 1C, the successive reduction in size of the mask 16 and the repeated dry etch processes are continued until the insulating material 14c and conductive material 12c is exposed, the mask 16 is removed, and a stair-step structure 10 remains. Word line contacts 20 are formed to extend through each respective insulating material 14 and electrically contact each conductive material 12, as shown in FIG. 1D. The top of each word line contact 20, as viewed in FIG. 1D, connects to a conductive word line (not shown). While FIGS. 1A through 1D illustrate using two anisotropic etches 18 to create three so-called “steps” of the stair-step structure 10, the acts of etching the insulating material 14, etching the conductive material 12, and trimming the mask 16 may be repeated to create more steps (and thus contact regions for word line contacts). Current conventional methods have been used to form more than eight contact regions (e.g., steps).\nAs the desired number of steps in the conventional stair-step structure increases, the margin of error associated with each act in the process of forming the steps correspondingly decreases when using the conventional method. For example, and as explained above, each iteration of the conventional method includes trimming the mask, etching the insulating material, and etching the conductive material. The desired number of steps is formed by repeating these acts as many times as the number of conductive materials in the stack. Each act of the conventional method has an associated etch control error because the size of each step is designed to fall within a particular range (e.g., tolerance) to allow enough room for a contact to be formed thereon while keeping the overall size of the stair-step structure small. Additionally, the relative locations of the steps are designed to fall within a range of locations in order to accurately form contacts thereon. As the number of iterations increases, any deviation from a target step width or location may be compounded because errors in one material are transferred to an underlying material. For a high number of steps in the stair-step structure, the margin of error to be achieved for the etch rate control may be less than one percent (1%). Small margins of error are difficult and costly to attain using conventional methods. Furthermore, because the mask is repeatedly trimmed, the method may start with a mask of high thickness, which may be difficult to repeatedly pattern and trim with the precision needed to have the necessary control over step width. Furthermore, the large amount of mask material is expensive and time-consuming to both form and remove.\nSpace savings in a memory device incorporating a vertical memory array may be accomplished by reducing the area that a stair-step structure covers. One method of reducing this area is described in U.S. Patent Application Publication No. 2009/0310415 to Jin et al., now U.S. Pat. No. 8,325,527, issued Dec. 4, 2012. Although some space is saved by aligning the word line contacts in the same direction as the bit lines, further improvements and reductions in cost in the manufacturing of such structures, as well as alternative methods of reducing the area covered by the stair-step structures, are desired. For example, the method described in Jin et al. uses a unique mask for each etch act to form the steps, which adds significant cost because of a high number of photolithographic reticles used to form the masks. Reductions in cost and improvements in controllability of manufacturing stair-step structures are, therefore, desired."}
{"text": "1. Field of the Invention\nThe present invention is related to dynamic random access memory (DRAM) devices. In particular, it relates to a system and method for improving a DRAM\"\"s ability to capture data correctly on all data paths (DQs).\n2. Description of Related Art\nWhen double data rate (DDR) is applied to a memory device, the data is input/output to/from the memory device on both edges (i.e., rising and falling) of an external system clock CLK. That is, the memory device (e.g., a DRAM) receives/outputs two bits of data per whole CLK cycle, called 2-bit (or 2n) prefetch. A DDR II memory device is similar to a DDR device but runs off an external clock CLK which is twice as fast (e.g., 200 MHz) as a DDR external clock. DDR II uses a 4-bit (or 4n) prefetch such that 4-bits are transferred in/out of a data path then handled as one 4-bit-wide piece of data inside the memory.\nDouble data rate memory is one key element for boosting memory device throughput to keep pace with the ever-increasing throughput performance of microprocessors. By doubling the memory bandwidth over current generation synchronous DRAM devices, DDR II can provide a cost-effective, high-performance main memory solution that does not require a significant development or manufacturing investment, while maintaining a cost structure consistent with synchronous DRAMs. DDR II memory devices and modules are well-suited for a broad range of applications, especially the workstation and server markets, where the high module density and device architecture can meet the performance and reliability demands of these products.\nDDR II memory technology has been defined and standardized by the Joint Electron Devices Engineering Council (JEDEC) as a next-generation memory solution. This technology is intended to facilitate adoption in a wide range of products, and to be offered in both device and module form from all major suppliers.\nThe DDR II standard proposes to have a minimum respective data setup and hold times of approximately 0.25 ns. Referring to FIG. 1, the data setup time is defined as the time the data is at the data input/output pin (DQ) before the next external system clock CLK edge arrives. The hold time is defined as the time the data stays valid after the clock edge. The entire window of data (setup and hold) is referred to as the xe2x80x9cdata eye.xe2x80x9d A data eye of approximately 1.5 ns is required in order for the DRAM to reliably capture the data (i.e., the DQ input pulse width (DIPW)). Moreover, in order for the DRAM to reliably capture the data, the clock CLK edge (either rising or falling) must be centered on the data eye. If the clock edge is not centered on the data eye, either one of the setup or hold times will be in danger of being of an insuffliciently short duration, which may prevent the DRAM from properly capturing the data.\nThere are many reasons why the clock edge may not be centered on the data eye. Some of those reasons include clock jitter, noise on the board, different lengths of data traces on the board, etc., which cause a clock skew to occur. The problem is exacerbated with higher clock frequencies since the setup and hold time is reduced in those cases, thus, leaving less room for error when attempting to locate the clock edge on the center of the data eye.\nOne overly complex solution to maintaining the clock edge near the center of the data eye that has been proposed is known as xe2x80x9cSyncLink.xe2x80x9d The SyncLink approach actually builds a mathematical model of the entire data eye and determines the exact location of the leading and trailing edges of the data. Once the system has calculated the location of the leading and trailing edges, the center of the data eye is calculated and the clock edge is located at the center of the data eye. While this method has proven to be accurate, it has also proven to be unnecessarily complex for most DDR II operating environments. This complexity results in greater chip complexity and cost. Thus, a simplified system and method are required to ensure that data which is input to a DRAM is properly timed relative to the external DRAM clock CLK such that the data is driven by a clock edge located on the center of the data eye and is, thereby, properly captured by the DRAM.\nThe present invention provides a simple system and method for controlling the rate with which data is input to a DRAM relative to the external DRAM clock CLK such that its clock edges are always centered on the data eye of input data to the DRAM, thus, ensuring that the data may always be properly captured by the DRAM.\nIn accordance with an exemplary embodiment of the invention, upon power-up or reset, during the existing initialization cycle time of the DRAM (i.e., 200 cycles under the DDR II standard), the write data capture calibration method of the invention is performed. A write command is issued by a memory controller and a simple repetitive string of data such as e.g., xe2x80x9c110011001100 . . . xe2x80x9d is sent by the memory controller to all the write data input paths (DQs) of the memory device, along with a bi-directional strobe (DQS) that is essentially a data clock aligned with data when it leaves the controller. The DRAM captures the data and performs a DQ to DQS delay adjustment to center the clock edge on the data eye for each DQ path. In the meantime, a delay lock loop (DLL) compares an internal data clock with the external clock to ensure they are both in phase. If they are out of phase, the DLL delays the internal clock until the two clocks are in phase. An auto-refresh automatically cancels the write data capture calibration method upon completion of the initial 200 cycles and the DRAM is returned to normal operation."}
{"text": "To meet the demand for wireless data traffic, which has increased since deployment of 4th-generation (4G) communication systems, efforts have been made to develop an improved 5th-generation (5G) or pre-5G communication system. Therefore, the 5G or pre-5G communication system is also called a ‘beyond 4G network’ or a ‘post long-term evolution (LTE) system’.\nIt is considered that the 5G communication system will be implemented in millimeter wave (mmWave) bands, e.g., 60 GHz bands, so as to accomplish higher data rates. To reduce propagation loss of radio waves and increase a transmission distance, a beam forming technique, a massive multiple-input multiple-output (MIMO) technique, a full dimensional MIMO (FD-MIMO) technique, an array antenna technique, an analog beam forming technique, and a large scale antenna technique are discussed in 5G communication systems.\nIn addition, in 5G communication systems, development for system network improvement is under way based on advanced small cells, cloud radio access networks (RANs), ultra-dense networks, a device-to-device (D2D) communication, a wireless backhaul, a moving network, a cooperative communication, coordinated multi-points (CoMP), reception-end interference cancellation, and the like.\nIn the 5G system, a hybrid frequency shift keying (FSK) and Feher's quadrature amplitude modulation (FQAM) and a sliding window superposition coding (SWSC) as an advanced coding modulation (ACM) scheme, and a filter bank multi carrier (FBMC) scheme, a non-orthogonal multiple access (NOMA) scheme, and a sparse code multiple access (SCMA) scheme as an advanced access technology have been developed.\nA communication system supporting a full-duplex scheme may increase system capacity doubles by performing a transmitting operation and a receiving operation on the same frequency at the same time.\nIn the communication system supporting the full-duplex scheme, there is a self-interference (SI) signal which occurs since a signal transmitted by a transmitting device is received in the transmitting device due to a characteristic of the full-duplex scheme.\nSo, various schemes for canceling the SI signal have been proposed in the communication system supporting the full-duplex scheme, and a typical one is a scheme for cancelling an SI signal which affects a receiving circuit of a transmitting device in a circuit domain.\nThe scheme for cancelling the SI signal in the circuit domain may be classified into a digital SI cancellation (SIC) scheme and an analog SIC scheme according to whether magnitude of the SI signal is within a digital dynamic range that the SI signal may be received in a digital domain. The digital SIC scheme denotes a scheme for cancelling an SI signal using a digital signal processing scheme, and the analog SIC scheme denotes a scheme for cancelling an SI signal using both an analog circuit and a digital signal processing scheme. It is general to use the digital SIC scheme and the analog SIC scheme at the same time for cancelling an SI signal in a system level.\nEach of the digital SIC scheme and the analog SIC scheme will be described below.\nFirstly, the analog SIC scheme will be described below.\nIn the analog SIC scheme, it will be assumed that a received SI signal includes a finite number of signals received after fixed delay time which a transmitting device already knows after the transmitting device transmits a signal. Under this assumption, the analog SIC scheme may adjust a gain for a transmission signal which is divided from an analog transmitting circuit included in the transmitting device, and cancel an SI signal from a signal received from the transmitting device by applying a circuit with fixed delay time. Here, the gain adjusted through the circuit may be acquired based on an interference characteristic estimated by the transmitting device.\nSecondly, the digital SIC scheme will be described below.\nThe digital SIC scheme detects a channel characteristic from difference between a transmission signal and a reception signal using a signal divided from a digital transmission signal of a transmitting device and a digital reception signal which is received using a digital signal processing scheme in the transmitting device, and cancels an SI signal from the digital reception signal by applying the channel characteristic in reverse.\nThe digital SIC scheme may cancel only an SI signal which is within a digital dynamic range. So, an SI signal which is not within the digital dynamic range may be cancelled only after an analog SIC operation is performed.\nSo, a performance of a communication system supporting a full-duplex scheme which uses the digital SIC scheme is determined according to a performance of the analog SIC scheme.\nThe analog SIC scheme is implemented with a scheme for previously predicting the number of SI signals which are received in a transmitting device after being reflected from the transmitting device, and adding an analog circuit which may cancel each SI signal.\nThe analog SIC scheme may not cancel an SI signal if a characteristic of the SI signal is different from a characteristic of an SI signal predicted by the analog circuit, or the number of SI signals is greater than the number of SI signals predicted by the analog circuit due to various reasons such as a situation that an SI signal is received from an outside of a transmitting device, and the like.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "The present invention relates to a coupler verification test circuit designed to assist personnel in installing a subscriber interface unit (SIU) on an optical fiber distribution network which does not have an operating head end, e.g. an operating office interface unit (OIU), during the installation procedure.\nOptical fiber is becoming widely preferred over electrical wires for use in telecommunication networks, e.g. telephony and video. Though optical fiber has far superior bandwidth than electrical wires and is also immune to electromagnetic and radio magnetic interference, problems exist with optical fiber networks during installation and configuration thereof. Specifically, when installing a subscriber to an optical fiber system during initial system installation, oftentimes the system does not have an operating head end. Accordingly it is difficult to impossible for a craftsman to readily determine if SIUs are being correctly installed onto the network during SIU installation. After the head end has been installed and is operating it can easily be determined if any or all the SIUs are functioning properly. However, at that time, though it is possible to determine which if any SIUs have been improperly installed, it is troublesome and expensive to detach personnel to correct the problems causing severe inconvenience."}
{"text": "A technique for imaging the distribution of molecular species in a living body is an important tool used in medical and biological research. Imaging of molecular species at the cellular level has been widely performed using a microscope and a molecular probe such as a molecular probe labeled with a fluorescence pigment or a chemiluminescence molecular probe. However, recently, there is a growing demand for devices for observing in vivo the distribution of molecular species of interest at the organ or whole-body level rather than the cellular level. For example, such an observation device allows the imaging of the distribution of target cancer cells labeled with a fluorescence probe in the body of a small living animal, such as a mouse, to monitor the growth of the target cancer cells over a fixed period of time, such as every day or every week. In a case where the growth of cancer cells in the body of an animal is monitored using a conventional device for cellular-level imaging, the animal is killed to stain or fluorescently-label cancer cells in a predetermined part of the body of the animal. In this case, the growth of cancer cells in the same individual cannot be monitored over a long period of time. For this reason, there is a demand for the development of a device capable of observing the distribution of molecular species in the body of a small living animal to obtain internal information about the body of the small animal.\nAs an exciting light-irradiating device for exciting fluorescence, one shown in FIG. 9 is known. As shown in FIG. 9, the exciting light-irradiating device has a filter wheel 8 and a multi-branched optical fiber bundle 16 to irradiate an object with light having a wavelength selected by the filter wheel 8 and the multi-branched optical fiber bundle 16. More specifically, light is emitted from a light source 2 such as a tungsten halogen lamp, collected by a lens 4 so as to enter an optical guide 6, and guided to a filter 10 mounted on the filter wheel 8 by the optical guide 6 so that only light passed through the filter 10 is guided to an entrance portion 16A of the multi-branched optical fiber bundle 16. Optical fibers constituting the multi-branched optical fiber bundle 16 are tied in a bundle at the entrance portion 16A, and are separated into four bundles at the position of a ring 16B provided on the way to a dark measurement chamber (not shown). The distal ends of these four optical fiber bundles are placed in the dark measurement chamber. The filter wheel 8 has a plurality of filters, and a desired excitation wavelength is selected by switching among these filters. Exciting light is guided by the multi-branched optical fiber bundle 16 to predetermined positions in the dark measurement chamber for measuring fluorescence. A device similar to the exciting light-irradiating device shown in FIG. 9 is also disclosed in U.S. Pat. No. 6,894,289."}
{"text": "The present application relates generally to an improved data processing apparatus and method and more specifically to mechanisms for performing data placement optimization using data context information collected during garbage collection operations.\nData context information represents object reference patterns during the execution of a program. Data context information provides valuable information for program understanding, performance analysis, and runtime optimizations. For object oriented applications, numerous objects are created, among which complicated reference patterns may occur. Therefore, data context information is very important for program understanding and optimizations.\nHowever, building data context information can be very expensive given the fact that (1) data accesses are highly frequent (a program can access a huge number of objects in a short period of time); and (2) the access patterns among objects are complex. Therefore, it is desirable to establish a program paradigm that can collect such information efficiently. Sacrificing some accuracy is affordable as long as enough information is available for the purpose of program understanding and optimization."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for producing a cap having a very small thickness to be fitted around a neck portion of a drink bottle such as a wine bottle or the like. Further, the present invention relates to a cap having a very small thickness to be fitted around a neck portion of a drink bottle such as a wine bottle or the like wherein the cap is produced using aluminum or aluminum alloy as a blank.\n2. Description of the Prior Art\nTo protect a neck portion of a drink bottle such as a wine bottle or the like from damage or injury, and moreover, maintain the neck portion in a clean state, a cap made of a metallic material is normally fitted around the neck portion of the bottle. Generally, the metallic cap is produced using a sheet of lead having both surfaces covered with tin foil. To plastically deform the sheet of lead to a contour corresponding to a cap product, a so-called spinning process has been heretofore employed wherein the sheet of lead is spun by manually actuating a specially designed tool while the sheet is rotated.\nWhen a metallic cap is produced by employing the foregoing spinning process, the thickness of a blank can not usually be reduced to 0.2 mm or less. Thus, with the conventional spinning process, there arises a drawback that material cost is increased because of the comparatively heavy thickness and the employment of an expensive metallic material like tin. Another drawback is that the surface of the cap produced by the spinning process can not exhibit a brilliant appearance. In addition, careless disposal after removal of the lead cap from the neck portion of a drink bottle may cause pollution. For this reason, it is anticipated that employment of lead for the cap will be prohibited in the future.\nGiven the circumstances described above, attention has recently turned to aluminum as a metallic material to be fitted around the neck portion of a bottle, because aluminum is a cheap metallic material, does not cause any public pollution after disposal and, moreover, can be reused by melting it.\nWith the conventional process, however, a cap having a very small thickness can be produced using a soft metallic material like lead but can not be produced when a comparatively hard metallic material like aluminum is employed as a raw material."}
{"text": "This invention relates generally to high voltage breakdown semiconductor devices and fabrication techniques therefor, and more particularly to a structure in which the lower junction termination of a multilayer semiconductor device formed in a wafer is extended to the top surface of the wafer, and a method of fabricating such structure.\nThere are many circuit applications for semiconductor devices which provide symmetrical blocking of applied voltages of different polarities and which exhibit high reverse breakdown voltages. In fabricating such device, it is necessary to control the geometries and characteristics of the device junctions. This can be accomplished by fabricating the device with a lateral, rather than vertical, structure. By bringing the lower junction termination of the device to the top surface of the wafer, better control of the symmetry and breakdown characteristics of the device can be achieved; however, this requires that the substrate of the device be electrically connected to the top surface of the wafer. While there are known ways to connect the substrate of a semiconductor device to its top surface, they generally have disadvantages which complicate the manufacture of semiconductor devices, such as the necessity of performing processing steps on each individual die. This is difficult and hence disadvantageous from a manufacturing standpoint.\nIt is desirable to provide a semiconductor device of relatively simple construction, having symmetrical blocking and high voltage breakdown characteristics, with its substrate being electrically connected to active layers on its upper surface, and a method of fabricating such a device which is suitable for large scale production. It is to this end that the present invention is directed."}
{"text": "Anti-glare (or antireflection) coatings on transparent substrates (e.g., glass) are important components for a large number of optical and optoelectronic devices, such as displays, lenses, and photovoltaic (PV) panels [1-10]. For instance, the unwanted optical reflection from the encapsulation glass layer of a PV panel could reduce the overall conversion efficiency of the solar device [11-13]. Anti-glare coatings are therefore widely applied on optical glass surfaces to reduce the reflection loss and increase the light transmission of the optical components [1,14]. Traditional quarter-wavelength antireflection coatings can effectively suppress optical reflection by satisfying the destructive interference conditions for the reflected light from the air/coating and the coating/substrate interfaces, thus decreasing the reflection and increasing the transmission of the substrate [14]. To satisfy the destructive interference conditions, the coating thickness needs to be close to one-fourth of the operating wavelength, while the refractive index of the coating (nc) needs to meet nc=√{square root over (nair×ns)}, when nair is the refractive index of air (1.0) and ns is the refractive index of the substrate [6,14]. For a typical glass substrate with a refractive index of 1.5, the anti-glare coating material needs to have a refractive index of ˜1.225. Low-refractive-index materials, such as MgF2 (with a refractive index of ˜1.37), are usually deposited on glass substrates by vacuum-based physical vapor deposition (PVD) technologies (e.g., sputtering) to achieve a precise control over the coating thickness [14-15]. Unfortunately, conventional PVD techniques suffer from high operating and equipment costs, limited material selection, low throughput, and small coating areas. These drawbacks particularly affect the applications where inexpensive anti-glare coatings on large-area glass substrates are needed, such as in solar industry.\nTo address the high costs and the low throughput issues of the vacuum-deposited anti-glare coatings, various simple solution processing technologies have been developed [1, 6, 16-25]. In many of these methods, nanoporous coatings with a large fraction of entrapped air and thus a low effective refractive index, which could satisfy the aforementioned ideal quarter-wavelength refractive index requirement, were extensively explored [1, 6, 16, 21, 26]. For example, nanoporous polymer coatings created by phase separation of spin-coated polymer blends, followed by selective removal of one component, have been demonstrated to show good anti-glare performance on glass substrates [1]. Multilayer silica nanoparticle coatings on glass substrates applied through common spin, dip, or roller coating techniques have already been commercialized for improving the efficiencies of PV panels (e.g., Honeywell's SOLARC RPV products) [27-29]. Electrostatics-assisted layer-by-layer (LBL) deposition of nanoparticles and polyelectrolyte multilayers is another popular approach in assembling anti-glare coatings on a variety of substrates [21, 30-31]. Monolayers of colloidal nanoparticles created by convective self-assembly [32-33], spin-coating [18, 29, 34-36], or Langmuir-Blodgett deposition [20, 37-38] have also been widely utilized as antireflection coatings on silicon and glass substrates. However, many of these existing wet-processing technologies involve multiple steps (e.g., LBL assembly) [19], are limited to single-sided coatings on planar substrates (approaches involving spin coating) [18], are not very reproducible over large areas [20], and/or are not inherently parallel for industry-scale manufacturing [32]. Thus, there is a need to overcome these deficiencies."}
{"text": "Many cardiac arrhythmias are caused by conduction defects that interfere with the propagation of normal electrical signals within the heart. The method adopted to treat arrhythmia is dependent on the nature and position of the underling conduction defect. Thus, electrophysiological mapping plays an important role in measuring the electrical activity of the heart. These techniques often require specialized equipment to locate the position of catheters in physical space and reconstructing the shape of the chamber from multiple site recordings. It would be desirable to provide 3D mapping without such equipment.\nState of the Art 3D mapping systems use magnetic fields, electrical fields or ultrasound to localize catheters. The main disadvantage of these systems is the prohibitive cost involved with the equipment and the need for both a conventional EP recording system and a separate 3D mapping/localization system. While manual positioning is not as accurate as current technologies, it is significantly more cost effective than conventional EP mapping systems and can be performed more rapidly.\nIt remains necessary to locate a target (active) site if an arrhythmia is to be terminated. A number of catheter locating systems are known in the art, but each introduces components and complexity to EP procedures. EP operators, however, are usually quite capable of piloting an EP catheter to a desired site within a patient's vasculature, particularly with fluoroscopic assistance. A difficulty remains, even if the location of the catheter is estimated based on fluorscopic guidance, in matching indwelling EP electrodes to sites on a cardiac model. This problem is all the more difficult when the model is rendered in 3D.\nIn part, the operator has data captured by a variety of systems. For example, electrogram channels monitor signals from indwelling electrodes, such as intracardiac electrodes and reference electrodes, and that information has to be coordinated with an anatomical (e.g., cardiac) model. Fluoroscopic images of the anatomy generally have no connection to other systems in the EP lab, and so piloting a catheter that lacks a locating system is done as a parallel, distinct part of the EP procedure. Cardiac mapping, therefore, has required great effort at a time when the operator's attention needs to focus on the patient or in labs where cost is an impediment and a highly trained technician is not available to operate a complex 3D mapping system.\nThe present invention addresses one or more of these problems."}
{"text": "Current architectures of projectors mainly include an illumination system, a light valve, and a projection lens. The illumination system is used for providing an illumination beam. The light valve is used for converting the illumination beam into an image beam. The projection lens is used for projecting the image beam onto a screen to form image pictures on the screen. The illumination system can generate illumination beams of different colors. A main principle is to excite phosphors on phosphor wheels through the light emitted from laser diodes having good luminous efficiency, and thereby to generate a desired pure color light source.\nIn light path architectures of projectors, a source of blue color light in the image pictures is a blue laser diode. In order to excite the phosphors on the phosphor wheels and achieve good efficiencies for conversion to different color lights, laser diodes that can emit short-wavelength blue light (i.e., blue light with a wavelength less than 445 nm) are used. But the short-wavelength blue light is a purplish blue light in the sense of human eyes. Therefore, in order to improve the quality of the image pictures, the purplish blue light must be adjusted to a blue light close to bluish. However, laser diodes have extremely narrow bandwidth with monochromaticity. Chromaticity coordinate values of the blue light, emitted by the laser diodes, in CIE1931 color space have been very close to critical values of the color space. Thus, the purplish blue light cannot be further purified through multicolor light filters to achieve desired chromaticity coordinate values in the color space.\nTherefore, how to improve above-mentioned problems is really a focus for relevant people in the field.\nThe information disclosed in this “BACKGROUND OF THE INVENTION” section is only for enhancement understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known to a person of ordinary skill in the art. Furthermore, the information disclosed in this “BACKGROUND OF THE INVENTION” section does not mean that one or more problems to be solved by one or more embodiments of the invention were acknowledged by a person of ordinary skill in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to a light adjusting apparatus that adjusts light by inserting/removing a light adjusting section into/from an optical path.\n2. Description of the Related Art\nImage pickup devices having an image pickup function are widely used in various fields, and among those fields, there is a field of small image pickup devices having a relatively small shape. Examples of such small image pickup devices include electronic endoscopes including a micro-video scope, optical microscopes provided with an image pickup function and portable devices provided with an image pickup function.\nSince downsizing is given priority in conventional small image pickup devices, a fixed focus lens, a fixed opening diaphragm, a fixed characteristic filter and the like are adopted as optical elements such as a lens, a diaphragm or an optical filter.\nIn contrast, high image quality is required also for such small image pickup devices in recent years, and there is a growing demand for adopting a focus lens, a variable diaphragm, and a variable characteristic filter as the aforementioned optical elements of light adjusting apparatuses, that is, functioning as a light adjusting apparatus that adjusts light.\nThus, many techniques are proposed which seek to downsize light adjusting apparatuses so as to be applicable to small image pickup devices.\nAs an example, Japanese Patent Application Laid-Open Publication No. 9-22042 describes an electromagnetic drive apparatus disposed around a taking lens which is provided with a yoke, a coil and a permanent magnet opposed to the yoke, the electromagnetic drive apparatus being configured to generate a magnetic force in the yoke by energizing the coil to cause the permanent magnet to rotate. By attaching, for example, a shutter blade rotatably and integrally to the permanent magnet as a light adjusting section, it is possible to switch between a state in which the shutter blade is positioned on an optical path and a state in which the shutter blade is retracted from the optical path."}
{"text": "1. Field of the Invention\nThe present invention relates to a phase-converter for converting a three phase motor to one phase operation.\n2. Description of the Prior Art\nIn many rural areas, the only electrical power available is single phase power from REA lines. In many cases, large electrical motors are required, for example, for operating pump jacks of oil pumping units. Large single phase electrical motors, however, are not available, and phase converters are required for converting large three phase motors to single phase operation. In the past, the control relay for switching the large starting capacitors out of the circuit after start up of the motor have been operated by A-C power from the motor which reacts to the line voltage. Problems have occurred in that there is much voltage fluctuation on the rural lines which causes the control relay to repetitively switch the large starting capacitors in and out of the circuit. In many cases this has caused the starting capacitors or motor to be destroyed."}
{"text": "1. Technical Field\nEmbodiments of the present invention relate to an adhesive composition for polarizing plates, a polarizing plate including the same, and an optical display including the same.\n2. Description of the Related Art\nA liquid crystal display includes a liquid crystal display panel and polarizing plates formed on both surfaces of the liquid crystal display panel. The polarizing plates may be stacked on the liquid crystal display panel via adhesive films for polarizing plates. Each of the adhesive films may be formed by coating an adhesive composition for polarizing plates onto one surface of the polarizing plate or a release film, followed by aging for a predetermined period of time.\nThe adhesive composition for polarizing plates is generally aged in an aging chamber for 3 to 7 days, or for 10 days or longer, as needed. As such, aging costs time and money. If the adhesive composition for polarizing plates is not subjected to sufficient aging, however, the properties of the adhesive film are poor (such as poor creep characteristics, peel strength, reliability, light leakage and reworkability), and the adhesive film therefore cannot be used in a polarizing plate. In addition, if the aging time is reduced, the adhesive composition can have poor pot-life, thereby causing poor coatability."}
{"text": "Sheets of material are often used in various industries and in a variety of ways. These materials can include paper, plastic, and other materials manufactured or processed in webs or sheets. As a particular example, long sheets of paper or other single layer or multi-layer products can be manufactured and collected in reels.\nThe “Z-structure” or cross-sectional structure of paper and other sheet materials is often a determining factor in numerous quality properties for the materials. For example, the distribution of voids between fibers in paper products typically affects bulk and opacity, and the distribution of fillers in paper products typically affects printing quality. Only fillers near the surface typically affect surface smoothness and ink permeability, and asymmetric filler distribution can cause the color of the surfaces of a paper sheet to differ.\nPapermakers are often interested in the distribution of fillers because there are several process adjustments available to influence it. As a result, accurate measurements of the Z-structure of a paper sheet could lead to the identification of filler distribution problems and allow timely adjustments to the paper-making process. Another issue may arise for intrinsically multi-layer sheets of material, such as those formed by splicing together multiple formed sheets or by coating a formed sheet with a polymer. In these cases, it is often desirable to know the thicknesses of individual layers or differences between exterior layers and interior layers."}
{"text": "1. Field of the Invention\nThe invention concerns keeping track of online research.\n2. Background Information\nOnline research has become a powerful tool for obtaining information on virtually any topic. Search engines provide an easy way to find information on Web pages. While finding the information may be relatively straightforward, capturing, saving, and organizing the information for later reference can be tedious. As a consequence, software tools have been developed that configure a computer to act as a Web-information manager, i.e., as a tool that performs in an automated fashion many of the tasks that people doing Web research had previously performed more manually."}
{"text": "The invention is generally related to mounting hardware, and more particularly to an extendable bracket for mounting window shades, curtain rods, and the like.\nThere are countless different kinds of bracket hardware for mounting or hanging window coverings such as shades, blinds, curtain rods, and the like. Many of the hardware designs are fabricated having a number of components including a mounting base for securing the hardware to a surface, and a support bracket for supporting the window covering components. Many of these hardware components are typically fabricated from metal and utilize several fastener components to complete a hardware assembly. Often, hardware brackets for hanging window coverings are not adjustable in any way.\nThe mounting base and support bracket are also typically fabricated separately and provided as a loose assembly to the consumer. The consumer then must install the mounting base, mount or attach the support bracket, and secure the mounting base and support bracket together using one or more additional fasteners. The installation process can therefore be quite cumbersome. In addition, the metal parts can be relatively costly to manufacture, both relative to labor, material expense, and tooling costs. In some cases, numerous parts are provided separately to the consumer, who must then assemble the parts prior to installation.\nMany bracket hardware designs are also adjustable in length to permit the installer to hang the window coverings at a desired distance from the mounting surface. Most of these bracket designs also utilize at least two separate components that telescope relative to one another and are thus slidable relative to one another to provide the length adjustability. Most of these designs have two or more metal parts that require multiple stamping operations to fabricate each part.\nSeveral known hardware brackets are length adjustable and utilize only a single bracket component per side of a window covering support rod. For example, U.S. Pat. No. 4,762,162 discloses a unitary hardware bracket for mounting window shades. The bracket is length adjustable by snapping or breaking off portions of the mounting end of the bracket. The wall mounting end can be shortened by breaking off one section of mounting holes and tangs. Another set of holes and tangs is left behind for mounting the bracket. However, if a consumer breaks off too much of the bracket during installation, rendering the bracket too short for the particular application, the only remedy is to purchase another bracket set.\nU.S. Pat. No. 2,752,991 discloses a window cornice mounting bracket assembly that also can be length adjusted at the wall mounting end by snapping off pieces of the bracket. Again, however, too much of the bracket can be accidentally or unintentionally broken off during installation.\nAnother problem with many current mounting hardware designs is that the brackets come only with a fixed length. A retailer and/or a manufacturer typically may offer a number of different bracket options. A series of similar brackets may be offered where each bracket in the series has a different, fixed length. In many cases, a retailer will offer, for example, four similar brackets having different lengths. The consumer must select the proper one, take it home, and install it. Many times, the selected bracket is either too short or too long for a particular job.\nTo illustrate, a consumer may be installing a curtain rod over a window that already has a blind installed. The curtain rod may need to extend further from the wall surface than the blind mounting brackets. The consumer therefore must either select the correct brackets the first time. If not, the consumer must return to the store, return the incorrect brackets, and re-purchase the correct brackets. This results in unnecessary extra effort for both the consumer and retailer. Alternatively, a consumer may choose to purchase a number of the different length brackets and use only the correct ones. The unused brackets likely will not be used, resulting in unnecessary expense and waste of resources.\nEither or both the manufacturer and retailer must also manufacture, stock, sip, unload, and keep track of each of the bracket options. Manufacturers and retailers often stock, store, display, and track product quantities and qualities according to computerized data, such as SKU numbers. Having four different brackets requires storage and shelf space for four different products and also requires tracking four different SKU numbers. This simply adds cost and complexity for both the manufacturer and the retailer."}
{"text": "An electromagnetic fuel injector comprises a cylindrical tubular body displaying a central feeding channel, which functions as a fuel conduit and ends with an injection nozzle regulated by an injection valve controlled by an electromagnetic actuator. The injection valve is provided with a needle, which is rigidly connected to a mobile keeper of the electromagnetic actuator in order to be displaced by the action of the electromagnetic actuator between a closed position and an open position of the injection nozzle against the bias of a spring which tends to hold the needle in the closed position. The valve seat is defined in a sealing element, which is shaped as a disc, lowerly and fluid-tightly closes the central channel of the support body and is crossed by the injection nozzle.\nPatent application EP1635055A1 describes an electromagnetic fuel injector in which a guiding element rises from the sealing element, such guiding element having a tubular shape, accommodating the needle therein in order to define a lower guide of the needle itself and displaying a smaller external diameter with respect to the internal diameter of the feeding channel of the supporting body so as to define an external annular channel through which pressurised fuel flows. Four through feeding holes, which lead towards the valve seat to allow the flow of pressurised fuel towards the valve seat itself, are obtained in the lower part of the guiding element. The needle ends with an essentially spherical shutter head, which is adapted to fluid-tightly rest against the valve seat and slidingly rests on an internal cylindrical surface of the guiding element so as to be guided in its movement. The injection nozzle is of the “multi-hole” type, i.e. it is defined by a plurality of through injection holes, which are obtained from a chamber formed downstream of the valve seat; in this way, the optimal geometries of the injection nozzle may be obtained for the various applications by appropriately orienting the single injection holes.\nExperimental tests have shown that the drive time-injected fuel quantity curve (i.e. the law linking the drive time to the quantity of injected fuel) of the electromagnetic injector described above is on the whole rather linear, but displays an initial step (i.e. displays a step increase for short drive times and therefore for small quantities of injected fuel); furthermore, the extent of such initial step is higher proportionally to the fuel feeding pressure.\nConsequently, the electromechanical injector described above may be used in a direct injection internal combustion Otto cycle engine (i.e. fed with petrol, LPG, methane or the like), in which the fuel feeding pressure is limited (lower than 200-250 bars) and the injector is not normally driven to inject small amounts of fuel). However, the electromagnetic injector described above cannot be used in a small direct injection internal combustion Diesel cycle engine (i.e. fed with Diesel fuel or the like), in which the feeding pressure of the fuel is rather high (up to 800-900 bars) and the injector is constantly driven so as to perform a series of pilot injectors before a main injection."}
{"text": "In the last few decades, the market for wireless communications devices has grown by orders of magnitude, fueled by the use of portable devices, and increased connectivity and data transfer between all manners of devices. Digital switching techniques have facilitated the large scale deployment of affordable, easy-to-use wireless communication networks. Furthermore, digital and radio frequency (RF) circuit fabrication improvements, as well as advances in circuit integration and other aspects have made wireless equipment smaller, cheaper, and more reliable. Wireless communication can operate in accordance with various standards such as IEEE 802.11x, Bluetooth, global system for mobile communications (GSM), code division multiple access (CDMA). As increased data throughput and other developments occur, updates and new standards are constantly being developed for adoption, such those associated with the third generation partnership project (3GPP).\nThe details of various embodiments of the methods and systems are set forth in the accompanying drawings and the description below."}
{"text": "The invention is based on an apparatus for the air-injection of liquid gas. An apparatus for the air-injection of liquid gas is already known but in which the intake tube pressure downstream of the throttle valve has an undesirable influence on the regulated mixture of liquid gas and air."}
{"text": "As a method for forming a high-quality semiconductor film, there is an epitaxial growth technique which forms a film on a substrate, using vapor phase growth. In a vapor phase growth apparatus using the epitaxial growth technique, a substrate is placed on a support portion in the vapor phase growth apparatus which is maintained at normal pressure or reduced pressure. Then, a reaction gas, which is a raw material, is supplied to the substrate while the substrate is being heated. For example, the thermal reaction of the reaction gas occurs on the surface of the substrate and an epitaxial single-crystal film is formed on the surface of the substrate.\nWhen a film is formed, the support portion is supported by a rotating body and the substrate is rotated by a bearing and a rotating mechanism which support a rotating shaft connected to the rotating body. When a film is formed, there is a concern that contaminants generated from the bearing or the rotating mechanism will be mixed with the film and will prevent a high-quality film from being formed."}
{"text": "1. Field of Invention\nThis invention relates to a system for manufacturing an instrument panel which is primarily used in the interior construction of an automobile.\n2. Description of the Relevant Art\nIn general, an instrument panel is constructed such that a foamed resin of polypropyrene or the like is placed between a hard core material (aggregate) of ABS resin or the like and a soft covering material such as polyvinyl chloride. In a conventional manufacturing process for such an instrument panel, beads of a foaming resin material are filled in a space between the core member and the covering material when formed, and are then heated and foamed. Due to this process, the prior art has various attendant problems in as much as a longer period of operation is required for making a foamed resin, and further, irritating gas is produced when the resin is foamed. In order to resolve such problems, proposals have been made by the present applicants in Japanese Patent Laid-Open No. 54-144, 478.\nIn the aforesaid known system a composite sheet having a semi-hard covering material and a foamed sheet adhered to each other in advance is cut to a desired length of a unit sheet and the unit sheet is guided to the heating station where it is heated to its softening point, thereafter the heated unit sheet is fed to the forming station where it is pressed against the core member set on the lower die in advance, through adhesive, and then finally it is pressed between the upper and lower dies. Further, in this prior art system of manufacturing an instrument panel, endless conveying chains are arranged to run through the feeding, heating and forming stations, respectively, and during this passage of the sheet panel, the sheet is spiked by some needles formed on the running chain the conveyed through each of the stations in sequence while being spiked. Due to this operational arrangement, the sheet should be conveyed through each of the stations with its positional relation to each of them being retained therein, and thus an independent processing of the sheet in each of the stations cannot be performed. Because the conveying chains pass through the heating station, they are always required to be cooled thus influencing the durability of the conveying chains. Further, an arrangement of the conveying chains in the forming station causes problems in that a sheet forming process is restricted and the like. With such arrangement as set forth above, two devices, one for drying the adhesive applied to the surface of the core member and the other for additionally heating the dried adhesive to melt it, are installed for processing the core member before the forming of the sheet material due to different temperature conditions therebetween.\nDue to this fact, an additional process for taking the core member from the drying device and transferring it to the melting device is required, and the area occupied by these devices is excessive, resulting in a decreased efficiency of the entire system. Further, because the heater, etc. should be arranged for each of the devices, the prior art system has the problem that a desirable thermal efficiency cannot be retained.\nFurther, in such a prior art system, the feeding station has a suction device for sucking the unit sheet cut to a desired length and a bogie device for transferring said unit sheet to the subsequent heating station at respective different locations. Due to this fact, an operation required for transferring the unit sheet to the bogie additionally requires time for transferring the sheet from the suction device to the bogie device, and further occupies an excessively large area the feeding station.\nFurther, in such a prior art system, problems arise in that the conveying chains themselves cannot be lifted up and down, and alternatively the upper and lower dies are required to be moved up and down and the lower die should be movable, resulting in a complex and large-sized device.\nIn the forming station, in order to set the core member on the lower die, an operator sets the core member while he enters a space between the dies and also in case of taking out the formed product of the instrument panel upon completion of the forming process, the product is pushed up by the ejector and the operator enters the space between the dies to take the product. Therefore, this is undesirable in view of workability and safety of the operator.\nThe present invention effectively overcomes the above problems found in a conventional type of an instrument panel manufacturing system as described above."}
{"text": "1. Field of the Invention\nThe present invention relates to a resist composition and a method of forming a resist pattern that uses the resist composition.\nPriority is claimed on Japanese Patent Application No. 2011-273759, filed Dec. 14, 2011, the content of which is incorporated herein by reference.\n2. Description of Related Art\nIn lithography techniques, for example, a resist film composed of a resist material is formed on a substrate, and the resist film is subjected to selective exposure followed by development, thereby forming a resist pattern having a predetermined shape on the resist film. A resist material in which the exposed portions of the resist film become soluble in a developing solution is called a positive-type, and a resist material in which the exposed portions become insoluble in a developing solution is called a negative-type.\nIn recent years, in the production of semiconductor elements and liquid crystal display elements, advances in lithography techniques have led to rapid progress in the field of pattern miniaturization. Typically, these pattern miniaturization techniques involve shortening the wavelength (and increasing the energy) of the exposure light source. Conventionally, ultraviolet radiation typified by g-line and i-line radiation has been used, but nowadays KrF excimer lasers and ArF excimer lasers are starting to be introduced in the mass production of semiconductor elements. Furthermore, research is also being conducted into lithography techniques that use an exposure light source having a shorter wavelength (and a higher energy level) than these excimer lasers, such as extreme ultraviolet radiation (EUV), electron beam (EB), and X-ray.\nResist materials for use with these types of exposure light sources require lithography properties such as a high resolution capable of reproducing patterns of minute dimensions, and a high level of sensitivity to these types of exposure light sources. As a resist material that satisfies these conditions, conventionally a chemically amplified resist composition has been used, which includes an acid generator component that generates acid upon exposure, and a base component that exhibits changed solubility in a developing solution under the action of acid.\nNumerous compounds have already been proposed for the acid generator component for chemically amplified resist compositions, including onium salt-based acid generators, oxime sulfonate-based acid generators, diazomethane-based acid generators, nitrobenzylsulfonate-based acid generators, iminosulfonate-based acid generators, and disulfone-based acid generators.\nResins (base resins) are typically used as the base components of chemically amplified resist compositions.\nFor example, in an alkali developing process where an alkali developing solution is used as a developing solution, a chemically amplified resist composition for forming a positive-type resist pattern typically contains an acid generator component and a resin component that exhibits increased solubility in an alkali developing solution under the action of acid. If the resist film formed using this resist composition is selectively exposed during formation of a resist pattern, then acid is generated from the acid generator component within the exposed portions, and the action of this acid causes an increase in the solubility of the resin component in an alkali developing solution, making the exposed portions soluble in the alkali developing solution. As a result, by performing alkali developing, the unexposed portions remain as a pattern, resulting in the formation of a positive-type pattern.\nAs the resin component, a resin that exhibits increased polarity under the action of acid is typically used. When the polarity of the resin is increased, the solubility in an alkali developing solution increases. On the other hand, when the polarity is increased, the solubility in an organic solvent decreases, and therefore if a solvent developing process that uses a developing solution containing an organic solvent (an organic developing solution) is employed instead of the alkali developing process, then within the exposed portions of the resist film, the solubility in the organic developing solution decreases relatively, meaning that during the solvent developing process, the unexposed portions of the resist film are dissolved in the organic developing solution and removed, whereas the exposed portions remain as a pattern, resulting in the formation of a negative-type resist pattern. This type of solvent developing process that results in the formation of a negative-type resist pattern is sometimes referred to as a negative-type developing process (for example, see Patent Document 1).\nCurrently, resins that contain structural units derived from (meth)acrylate esters within the main chain (acrylic resins) are widely used as base resins for chemically amplified resist compositions designed for use in ArF excimer laser lithography or the like, as they exhibit excellent transparency in the vicinity of 193 nm (for example, see Patent Document 2). Here, the term “(meth)acrylate ester” is a generic term that includes either or both of the acrylate ester having a hydrogen atom bonded to the α-position and the methacrylate ester having a methyl group bonded to the α-position. The term “(meth)acrylate” is a generic term that includes either or both of the acrylate having a hydrogen atom bonded to the α-position and the methacrylate having a methyl group bonded to the α-position. The term “(meth)acrylic acid” is a generic term that includes either or both of acrylic acid having a hydrogen atom bonded to the α-position and methacrylic acid having a methyl group bonded to the α-position.\nIn order to improve the lithography properties and the like, the base resin typically includes a plurality of structural units. For example, in the case of an aforementioned resin component that exhibits increased polarity under the action of acid, typically, a base resin is used that contains a structural unit having an acid-decomposable group that decomposes under the action of the acid generated from the acid generator component, resulting in increased polarity, and also contains a structural unit having a polar group such as a hydroxyl group, and a structural unit having a lactone structure and the like.\nOne known technique for further improving the resolution is a lithography technique known as liquid immersion lithography (hereafter also referred to as “immersion exposure”), in which exposure (immersion exposure) is conducted in a state where the region between the objective lens of the exposure apparatus and the sample is filled with a liquid (an immersion medium) that has a larger refractive index than the refractive index of air.\nBy using immersion exposure, it is considered that higher resolutions equivalent to those obtained using a shorter wavelength light source or a higher NA lens can be achieved using the same exposure light source wavelength, with no reduction in the depth of focus. Furthermore, immersion exposure can be conducted using existing exposure apparatus. As a result, it is expected that immersion exposure will enable the formation of resist patterns of higher resolution and superior depth of focus at lower costs, and in the production of semiconductor elements, which requires enormous capital investment, immersion exposure is attracting considerable attention as a method that offers significant potential to the semiconductor industry, both in terms of cost and in terms of lithography properties such as resolution.\nImmersion lithography is effective in forming patterns having all manner of shapes. Further, immersion exposure is capable of being used in combination with super-resolution techniques such as phase shift methods and modified illumination methods that are currently under investigation. Currently, techniques using an ArF excimer laser as the exposure source are the most actively researched immersion exposure techniques. Further, water is mainly being investigated as the immersion medium.\nIn recent years, the addition of a photoreactive quencher to a chemically amplified resist composition has also been proposed (for example, see Patent Documents 3 and 4). A photoreactive quencher is a salt formed from an anion and a cation, which, prior to exposure, has a quenching action that traps acid generated from the acid generator or the like via an ion exchange reaction, but which decomposes upon exposure, resulting in a loss of the quenching action. Accordingly, when a resist film formed using a chemically amplified resist composition containing such a photoreactive quencher is subjected to exposure, in the exposed portions, the photoreactive quencher loses its basicity relative to the acid generated from the acid generator or the like, whereas in the unexposed portions, the photoreactive quencher traps the acid, thereby suppressing the diffusion of acid from the exposed portions into the unexposed portions, resulting in improved lithography properties."}
{"text": "1. Field\nSubject matter disclosed herein relates to glucose sensor signal stability analysis including, by way of example but not limitation, analyzing a reliability of a glucose sensor signal by attempting to detect a change in responsiveness of the sensor signal.\n2. Information\nThe pancreas of a normal healthy person produces and releases insulin into the blood stream in response to elevated blood plasma glucose levels. Beta cells (β-cells), which reside in the pancreas, produce and secrete insulin into the blood stream as it is needed. If β-cells become incapacitated or die, which is a condition known as Type I diabetes mellitus (or in some cases, if β-cells produce insufficient quantities of insulin, a condition known as Type II diabetes), then insulin may be provided to a body from another source to maintain life or health.\nTraditionally, because insulin cannot be taken orally, insulin has been injected with a syringe. More recently, the use of infusion pump therapy has been increasing in a number of medical situations, including for delivering insulin to diabetic individuals. For example, external infusion pumps may be worn on a belt, in a pocket, or the like, and they can deliver insulin into a body via an infusion tube with a percutaneous needle or a cannula placed in subcutaneous tissue.\nAs of 1995, less than 5% of the Type I diabetic individuals in the United States were using infusion pump therapy. Over time, greater than 7% of the more than 900,000 Type I diabetic individuals in the U.S. began using infusion pump therapy. The percentage of Type I diabetic individuals that use an infusion pump is now growing at a rate of over 2% each year. Moreover, the number of Type II diabetic individuals is growing at 3% or more per year, and increasing numbers of insulin-using Type II diabetic individuals are also adopting infusion pumps. Physicians have recognized that continuous infusion can provide greater control of a diabetic individual's condition, so they are increasingly prescribing it for patients.\nA closed-loop infusion pump system may include an infusion pump that is automatically and/or semi-automatically controlled to infuse insulin into a patient. The infusion of insulin may be controlled to occur at times and/or in amounts that are based, for example, on blood glucose measurements obtained from an embedded blood-glucose sensor, e.g., in real-time. Closed-loop infusion pump systems may also employ the delivery of glucagon, in addition to the delivery of insulin, for controlling blood-glucose and/or insulin levels of a patient (e.g., in a hypoglycemic context). Glucagon delivery may also be based, for example, on blood glucose measurements that are obtained from an embedded blood-glucose sensor, e.g., in real-time."}
{"text": "This invention relates, in general, to transmitter optoelectronic integrated circuits, and more particularly to a distributed drive, vertically integrated, transmitter optoelectronic integrated circuit for mating with a plastic optical fiber.\nOptoelectronic transmission systems have emerged as a prominent technology in a variety of disciplines including: automotive, computer, medical, and communications. Typically, an optoelectronic system comprises a transmitting optoelectronic integrated circuit, or transmission source, coupled to a receiving integrated circuit via an optical fiber. Two key indices of an optoelectronic signal transmission system are the coupling efficiency between the transmission source and the fiber optic cable, and the power level of the signal generated by the transmission source.\nCoupling efficiency is a figure of merit that indicates how much of the optical signal generated by the transmission source is conducted by the optical fiber. The theoretical maximum coupling efficiency is a function of the dimensions of both the optical fiber and the transmission source as shown by Hudson in his paper \"Calculation of the Maximum Optical Coupling Efficiency into Multimode Optical Waveguides\" (Applied Optics Vol. 13, No. 5, May 1974). In particular, the maximum coupling efficiency is determined by the diameter and numerical aperture of the optical fiber as well as the diameter of the emission source. Hence, the optical fiber is an important component in the optoelectronic transmission system.\nGenerally, optical fibers are cylindrically shaped with an inner core surrounded by an outer core, commonly referred to as a cladding layer. Two optical fiber parameters which strongly influence maximum coupling efficiency are the inner core diameter and the numerical aperture. In the past, the preferred material for the optical fiber has been glass; a material in which both the inner core diameter and numerical aperture are relatively small. Further, the use of glass optical fibers requires that the transmitting and receiving portions of these systems employ expensive packaging materials to ensure adequate coupling between these two components.\nAccording to the mathematical relationship derived by Hudson, selecting an inner core diameter and a numerical aperture of the optical fiber constrains the diameter of the transmission source for a selected maximum coupling efficiency. Further, to obtain an acceptable maximum coupling efficiency, the relationship between the inner core diameter and the numerical aperture of the optical fiber limits the transmission source diameter to be a small percentage of the inner core diameter. As an example, for a typical glass fiber with an inner core diameter of 50 micrometers and an numerical aperture of 0.21, the transmission source diameter is limited to 20 percent of the optical fiber inner core diameter for the theoretical maximum coupling efficiency.\nThe primary disadvantage of a small transmission source diameter is that the injection current density in the transmitting device must be relatively high to achieve an acceptable minimum coupled power from the transmitting optoelectronic integrated circuit to the optical fiber. Moreover, thermal properties of both transmitter optoelectronic semiconductor devices, and the lower cost plastic packages used for encapsulating the devices, limit the maximum injection current. Hence, theoretical maximum coupled power has been constrained by physical limitations of the transmitting semiconductor device in addition to those posed by the optical fiber cable.\nFurther, the use of glass optical fibers has limited the topography of the circuitry associated with transmitter optoelectronic integrated circuits such that current must be collected at contact regions of a transmission source transistor, and transported through metal interconnects. Ultimately, the current must be redistributed by an optimized ohmic contact pattern on an optical emission device. Associated with this current distribution scheme are parasitic resistances and capacitances that degrade the performance of the transmitter optoelectronic integrated circuit. And as discussed previously, the use of glass optical fibers requires very sophisticated and expensive optoelectronic circuit packaging material to achieve the theoretical maximum coupling efficiency.\nMore recently plastic optical fibers have gained widespread acceptance as an alternative to glass optical fibers. Plastic optical fibers have both a larger inner core diameter and numerical aperture than do their glass counterparts. Hence, for a similar maximum coupling efficiency, the transmission source diameter for a plastic optical fiber may be greater than that of a glass optical fiber. As an example, to achieve the theoretical maximum coupling efficiency, a typical plastic optical fiber with an inner core diameter of 1000 micrometers and a numerical aperture of 0.47 the transmission source diameter can be up to 50 percent of the diameter of the plastic optical fiber. Hence, one advantage derived by using plastic optical fibers is that the larger allowable transmission source diameter provides the option for a lower injection current density. Further, plastic optical fibers make the use of lower cost plastic packaging for the transmission optoelectronic devices feasible. Unfortunately, the present methods for fabricating the transmission optoelectronic devices still renders the use of high current densities in localized portions of the optoelectronic devices. Further, the parasitic resistances and capacitances associated with current collection, transport, and redistribution in the optoelectronic devices still exist.\nAccordingly, it would be beneficial to have a transmitter optoelectronic integrated circuit capable of achieving high output current levels with relatively low current densities; while simultaneously taking advantage of the increased transmission source area afforded by using plastic optical fibers to optimize the optoelectronic integrated circuit layout wherein the circuit does not suffer from the performance degradations caused by parasitic resistances and capacitances associated with current collection, transport, and redistribution."}
{"text": "There is a need to identify objects and people (“objects”) automatically. Automatically determining identity eliminates human error, makes the identification process faster, and allows direct interfacing with computing systems and related applications. The most common types of automatic identification in use today include barcodes, RFID labels, and magnetic cards. Automatic identification is used for applications such as identifying items in inventory or for checkout, providing entry into secured areas, logging into a computer system, facilitating rental card returns, positively identifying patients in hospitals, ensuring the proper medication is being delivered, and numerous other applications. Unfortunately the systems used today have many drawbacks and inadequacies. Conventional systems typically require the user to hold a reading device (such as a barcode scanner), or to hold an identification device (such as an ID card) thus tying up the user's hands. In addition, these systems only work when a barcode, RFID tag, or magnetic card is properly scanned. Most scanners only work a very short distance, require that the barcode, RFID tag, or magnetic card be correctly oriented to the reader, and frequently do not read correctly on the first scan, making them inadequate for many applications. Real time location systems (“RTLS”) can also provide automatic identity, but are large and complex systems that are expensive and difficult to install and maintain. In addition most location systems do not have the spatial or temporal accuracy required for most automatic identification applications.\nThe drawbacks and inadequacies of current systems create significant compromises in many applications. For example ID cards containing a barcode, magnetic strip, or RFID chip are commonly used to unlock a door for access to restricted areas. People desiring to enter the area must scan the card at a reader near the door. The person must retrieve the ID card from a pocket, wallet, or purse, and orient it correctly to the reader. If the card is not oriented correctly or not read properly by the reader, the person must repeat the operation. This process can be time consuming and frustrating to the user.\nSimilar ID card systems are used to authorize access to computer systems. Frequently these systems are used together with a password to provide positive user identification and access to the system. These systems have many of the same problems as entry access systems. People desiring to access the computer must scan the card at a reader. The person must retrieve the card from a pocket, wallet, or purse, and orient it correctly to the reader. If the card is not oriented correctly or not read properly by the reader, the person must repeat the operation. This process can be time consuming and frustrating to the user. Frequently the user will set the ID card down on the work surface next to the computer and forget to pick it up when they are finished with their task, which can lead to a lost card and/or a security risk. Another problem with these systems is that people frequently forget to logoff when they have completed their tasks and simply walk away from the computer. When this happens, an unauthorized person can then walk up to the computer and have unauthorized access. In an attempt to prevent this problem from occurring some computer systems automatically log a user off if there is isn't keyboard activity for some period of time. If the timeout period is too short then users gets logged off inadvertently. A user may turn to talk briefly to a colleague and turn back to the computer only to find that they have been logged off. If the timeout period is too long and the user forgets to logoff then there is a period of venerability for unauthorized access.\nThe need for an improved computer authorization system is particularly evident in hospitals where mobile care providers access many different computer systems on a frequent basis and are faced with strict HIPPA rules for patient confidentiality. Hospitals can also benefit from other applications of automatic identification systems. Another example is positive patient identification. Clinical personnel need to positively identify patients prior to administering medication or performing procedures. Hospitals have tried to utilize barcodes and RFID tags located on patient's wrist bands for this purpose. This requires that the clinician use a reading device which means the clinician must be located next to the reading device, search for a portable device, or ensure that he has a portable reader on him at all times. This creates issues in spending time searching for a reader, or the hassle of carrying a reader. There is also the issue of limited battery life in portable readers, where even if the clinician has a device available it might not be usable due to a low battery. In addition it can be difficult to scan the barcode or RFID tag since it is located on the patient's wrist and can be in a variety of positions making it difficult to properly align the reader or even get the reader near the ID band.\nA touch identification system could also enable new applications such as automatically determining whether personnel are complying with hand washing requirements. In many jobs such as food-workers and hospital-workers there are requirements for hand-washing to ensure public safety. Unfortunately these policies are frequently not followed. To ensure compliance it would be desirable to have a system that automatically determined that an employee activated the soap dispenser and logged the event.\nTechnologies and products have been developed to provide communication of information through or on the body of a user, but these technologies and products have not been broadly adopted due to their significant drawbacks and inadequacies. Many of these technologies and products require very specific placement of devices on the body of the user and/or require multiple contact points—some including ground connections—making them impractical to use. Most of these technologies and products do not address the issue of communication signals inadvertently radiating or coupling from the devices or the user, which could result in signals being wrongly communicated creating security issues with the system. There is also the issue of unreliable performance due to users having different body characteristics such as body mass and skin impedance and environmental issues causing variable stray coupling of the user to objects in the environment."}
{"text": "Citalopram and its pharmaceutically acceptable acid addition salts, such as its hydrogen bromide salt (Formula (I)) as described in U.S. Pat. No. 4,136,193, are anti-depressant drugs with few side effects.\n\nVarious processes for the preparation of citalopram have been described in the prior art. For example, U.S. Pat. No. 4,136,193 describes the alkylation of cyanophthalane with 3-dimethylaminopropylchloride, using sodium hydride as a base in a dimethylsulphoxide (DMSO) medium (Scheme-1). The reaction mixture is poured into ice water and extracted with ether. Then, after standard acid-base work-up, crude citalopram base is isolated as an oil. The isolated oil is purified by high vacuum distillation, (0.03 mm at 175-180° C.) and then converted to the hydrobromide salt.\n\nAnother process, as described in WO-A-98/19511, is the alkylation of cyanophthalane with 3-dimethylaminopropylchloride in the presence of a strong base (such as n-butyl lithium) and diisopropylamine in a dimethoxyethane medium at −50° C. (Scheme-2). After completion of the reaction, the reaction mixture is poured into ice water and extracted with toluene. After standard acid-base work-up using toluene as solvent, citalopram base is isolated as an oil. The oily base is then converted to acid addition salts such as citalopram hydrobromide and hydrochloride by conventional methods.\n\nA drawback of the above two processes is that the citalopram base is isolated as an oil. Purification of citalopram oily base is carried out by high vacuum distillation (0.03 mm) at 175-181° C. Achieving such a very high vacuum at plant level is difficult and hence the process described is not easily transferable to the commercial scale. Apart from these constraints, citalopram base having a cyano group at the 5th position of the bicyclic ring system may decompose during high vacuum distillation at high temperature to form citalopram carboxamide as an impurity, resulting in poor quality and yield.\nIn yet another process, citalopram is made as described in U.S. Pat. No. 4,650,884. The process involves the successive Grignard reaction of 5-cyanophthalane with 4-fluorophenylmagnesiumbromide and N,N-dimethylaminopropylmagnesiumchloride and the cyclization of the resulting diol to obtain citalopram base as an oil (Scheme-3). The oily base is converted to the hydrobromide salt using anhydrous hydrogen bromide gas in acetone medium. Due to the poor quality of the oily base, repeated crystallization of the hydrobromide was necessary to obtain a pharmaceutically acceptable quality of citalopram hydrobromide.\n\nUsing the same strategy, citalopram base may be made as described in CA-A-1,339,452 [equivalent to U.S. Pat. No. 4,943,590, EP-A-347066 & GB-A-8,814,057 Scheme-3]. The racemic diol base, as described in the patent, is dissolved in dichloromethane. Triethylamine and methane sulphonyl chloride are added over a period of 1 hour. The reaction mixture is then washed twice with 0.1M sodium hydroxide solution, the organic phase is separated and dried over anhydrous magnesium sulphate followed by solvent concentration under reduced pressure, and the citalopram base is isolated as a crystalline solid. For the first time, citalopram base has been reported as a crystalline solid. However, no physical data are described here [EP-A-347066]. The solid base is converted to the hydrobromide salt using anhydrous hydrogen bromide gas in acetone medium.\nDE-A-20 007 303 discloses yet another process for the isolation of citalopram base, as a solid, from citalopram hydrobromide. However, the process for making citalopram hydrobromide is not described. The isolated crystalline base is then converted into the desired salt. According to the process described here, pure citalopram hydrobromide is dissolved in 5 volumes of water and the pH is adjusted to about 10 with 6N sodium hydroxide. Citalopram base is then extracted into a non-polar organic solvent, such as toluene. The toluene is distilled off under reduced pressure and the resulting residue is triturated with n-heptane to precipitate citalopram base as a solid. The solid is then filtered to produce crystalline citalopram base.\nThe main disadvantage of this process is that citalopram hydrobromide has to be made first from the crude oily base isolated from the prior art process and then converted back to solid citalopram base. Then, the solid base is again converted to its corresponding hydrobromide salt. Thus, the process involves multiple operations which are tedious and time consuming. In addition, prolonged heating of citalopram base may increase the carboxamide impurity, resulting in poor quality citalopram."}
{"text": "The present invention is directed to a communication cable with a cable core provided in the center surrounded by an outside cladding as well as strain relief elements applied in the region of the outside cladding.\nAllowed U.S. patent application Ser. No. 08/646,322 which issued as U.S. Pat. No. 5,706,381, whose disclosure is incorporated by reference and which is a national phase application based on WO-A1 95/13556, discloses a cable having strain relief elements adjacent the outside surface of the cable. The tensile elements are arranged in the cladding region and lie in a symmetrical plane that simultaneously forms the bending plane for the cable. The tensile elements lie diametrically opposite one another and thus are offset by 1800 from one another. The light waveguides are arranged in the center of the cable in a single layer or in a multi-layer envelope.\nIn many instances, it is desirable to construct cables in a type of hybrid structure, for example to arrange a plurality of transmission elements that are at least partially independent of one another in the cable. For example, this is the case where both electrical as well as optical transmission elements are employed. However, applications with two or more optical transmission elements are also conceivable in this context."}
{"text": "This disclosure relates to a method of making an antistatic agent.\nThermoplastics are useful in the manufacture of articles and components for a wide range of applications, from automotive parts to electronic appliances. Because of their broad use, particularly in electronic applications, it is desirable to provide thermoplastic resins with antistatic agents. Many polymers or blends of polymers are relatively non-conductive, which can lead to static charge build-up during processing and use of the polymer. Charged molded parts, for example, may attract small dust particles, and may thus interfere with a smooth surface appearance, for example by causing a decrease in the transparency of the article. In addition, the electrostatic charge may be a serious obstacle in the production process of such polymers.\nAnti-static agents are materials that are added to polymers to reduce their tendency to acquire an electrostatic charge, or, when a charge is present, to promote the dissipation of such a charge. Organic anti-static agents are usually hydrophilic or ionic in nature. When present on the surface of polymeric materials, they facilitate the transfer of electrons and thus eliminate the build up of a static charge. Anti-static agents have also been added to the polymer composition before further processing into articles, and may thus be referred to as “internally applied.” Useful anti-static agents applied in this manner are thermally stable and able to migrate to the surface during processing.\nA large number of anti-static agents having surfactants as their main constituent have been considered and tried. Many suffer from one or more drawbacks, such as lack of compatibility with the polymer (which interferes with uniform dispersibility), poor heat stability, and/or poor antistatic characteristics. Poor heat resistance in particular can adversely affect the optical properties of engineering thermoplastics such as aromatic polycarbonates.\nParticular phosphonium salts of certain sulfonic acids, however, have been shown to be useful antistatic agents. U.S. Pat. No. 4,943,380 discloses reducing the static charge on polycarbonate resins with an anti-static composition containing 90-99.9 weight % of polycarbonate and 0.1-10 weight % of a heat resistant phosphonium sulfonate having the general formula:\nwherein R is a straight or branched chain alkyl group having 1 to 18 carbon atoms; R1, R2 and R3 are the same, each being an aliphatic hydrocarbon having 1 to 8 carbon atoms or an aromatic hydrocarbon group having 6 to 12 carbon atoms; and R4 is a hydrocarbon group having 1 to 18 carbon atoms.\nU.S. Pat. No. 6,194,497 discloses antistatic resin compositions, particularly transparent resin compositions, comprising a thermoplastic polymer and a halogenated medium- or short-chain alkylsulfonic acid salt of a tetrasubstituted phosphonium cation. The antistatic agent described therein is prepared by ion exchange of a potassium haloalkylsulfonate to produce the corresponding acid. The haloalkylsulfonic acid is then reacted with tetrabutylphosphonium hydroxide to product the antistatic agent.\nAn advantage of this synthesis is that use of an ion exchange step during synthesis results in a product that is very pure, i.e., contains little to no halogenated compounds that may ultimately lead to degradation of resins such as polycarbonates. However, while suitable for its intended purposes, this particular synthesis also has a number of drawbacks. For example, use of an ion exchange step increases the expense of the process, and may lead to the production of waste requiring disposal procedures. The synthesis also uses the potassium salt as a starting product, which is prepared from the corresponding sulfonylfluoride. Since the solubility of potassium peralkylsulfonates is relatively low, e.g., on the order of 5% at 20° C., a water/ethanol mixture is needed in the ion exchange. The flammability of ethanol requires the implementation of significant safety precautions during the synthesis. In addition, selecting the appropriate water/ethanol ratio is also important. An excess of alcohol may render the final product soluble in the reaction solvent, such that isolation of the product may require a further extraction step.\nThere accordingly remains a demand in the art for more efficient processes, particularly one-step processes, for making phosphonium sulfonate antistatic agents, as well as thermoplastic resin compositions that incorporate these antistatic agents. It would further be desirable for such processes to produce the antistatic agent in good yields without having a detrimental effect on the safety of the process and/or the purity of the product."}
{"text": "A social network is a set of people (or organizations or other social entities) connected by a set of social relationships, such as friendship, co-working or information exchange relationship. There has been a recent unprecedented increase in the use of Online Social Networks (OSNs) to expand our social life, such as finding others of a common interest, discussing and sharing information in forums, and exchanging photos and personal news. The OSNs have become a large-scale distributed system providing services to hundreds of millions of users and delivering messages at very high rate.\nBesides handling traditional client-to-server requests, OSNs also need to handle highly interconnected data due to the strong community structure and human relationships among their end users, which often results in complex data sharing among users. Given the tremendous user population and frequent data access by these users, effective resource planning and provisioning strategies are of extreme importance to the performance and revenue of an OSN. In particular, selecting the most suitable locations to deploy server farms is one of the key steps in such resource management.\nThe development of placement strategies for resources including servers and bandwidth affects the performance of any online web services. An appropriate allocation of resources benefits content providers by reducing latency for their clients and balancing the bandwidth consumption. The goal is to provide content distribution to clients with good Quality of Service (QoS) while retaining efficient and balanced resource consumption of the underlying network infrastructure. Thus, existing server placement proposals mainly focus on minimizing the average latency between the server and the users, given the nature of client/server communication patterns in traditional web services.\nMany proposals on the server placement problem rely on extracting clients' requests from history traces collected on the web servers, and then searching for the best placement given the particular client and load distribution. While these proposals might be plausible in improving performance of existing OSN services, these proposals are less helpful to new OSNs that are starting afresh. Thus, a problem arises in that it is difficult for a new born Internet application service to make a decision on where to deploy its servers.\nMoreover, much existing work on server placement casts the problem as an integer linear program where a binary decision variable bij is used to denote if user i is assigned to server j; and the total number of selected servers should be a predetermined input M. One of the best known approximation algorithms for this problem (presented by M. Charikar and S. Guha in “Improved combinatorial algorithms for the facility location and k-median problems” in Proceedings of the 40th Annual Symposium on Foundations of Computer Science, 1999) achieves a very large time complexity of O((N+P)3), where N is the number of servers and P is the number of users. To make the problem more manageable in reality, a number of approximation and heuristics have been proposed such as the use of a greedy algorithm (by L. Qiu, V. N. Padmanabhan, and G. M. Voelker, in “On the Placement of Web Server Replicas,” in Proc. of IEEE INFOCOM 2001, 1587-1596). However, these approaches are mainly based on theoretical analysis and are only validated using simulation in very small graphs. Further, such methods have fundamental issues in scaling to a large value of N. Accordingly, there is a need for methods to efficiently and flexibly determine, for existing or new OSNs and any value of N, where to place servers."}
{"text": "Most refrigerators intended for household use include an ice maker and an ice bin, which generally both stores the ice and provides access to the ice. Depending on the configuration of the refrigerator and/or the placement of the ice maker, accessing or reaching the ice may be difficult. Also, the amount of ice that can be produced and stored at one time is limited by the size of the ice bin.\nA number of refrigerators include an ice dispenser coupled to the ice bin that dispenses ice from the ice bin through a refrigerator door. Typically, the user operates the ice dispenser by pushing a drinking glass against a paddle or other lever. When the lever is depressed, ice is released directly from the ice bin into the glass. While this may simplify accessing the ice, retrieval of the ice is limited to the dispenser's speed (and the size of the glass). As a result, removal of large amounts of ice using the dispenser is difficult and time-consuming."}
{"text": "In data-type image processing, high demands are made on a very fast performance of the image-associated operations. A way of achieving a high processing speed is to use units with an integrated unit solution comprising both a camera in the form of a photodiode matrix and an image processor. Such an arrangement is known from, e.g., SE, B, 431 145 which presents an image processor which contains a device able to indicate the image dots belonging to a connected object, and produces a binary signal corresponding to the object when the relevant image dot, on the one hand, meets a condition specific to the object, and, on the other hand, is indicated as entering the object, or when the closest adjacent image dots are distinguished. By means of this known technique the need is reduced for a serial output of image data, and a substantially reduced amount of information is obtained which, with less time consumption than for the known image processing systems, can be output for further image processing.\nAnother example of corresponding parallel signal processing is described in Swedish Patent No. 9001556-1. The intention here is to be able to carry out computations on radar signals in addition to image processing.\nA drawback with the known technique is that the known processors cannot perform so-called arithmetic global operations with the same efficiency, for example determine the centre of median for every object in a read-in vector. In the known device, there takes place with this type of operation a serial outputting of image data which, in per se time terms, reduces the advantages of the fast image processors."}
{"text": "Agents, such as therapeutic agents, can be delivered systemically, for example, by injection through the vascular system or oral ingestion, or they can be applied directly to a site where treatment is desired. In some cases, particles are used to deliver a therapeutic agent to a target site."}
{"text": "Patient vital sign monitoring may include measurements of blood oxygen, blood pressure, respiratory gas, and EKG among other parameters. Each of these physiological parameters typically requires a sensor in contact with a patient and a cable connecting the sensor to a monitoring device. For example, FIGS. 1-2 illustrate a conventional pulse oximetry system 100 used for the measurement of blood oxygen. As shown in FIG. 1, a pulse oximetry system has a sensor 110, a patient cable 140 and a monitor 160. The sensor 110 is typically attached to a finger 10 as shown. The sensor 110 has a plug 118 that inserts into a patient cable socket 142. The monitor 160 has a socket 162 that accepts a patient cable plug 144. The patient cable 140 transmits an LED drive signal 252 (FIG. 2) from the monitor 160 to the sensor 110 and a resulting detector signal 254 (FIG. 2) from the sensor 110 to the monitor 160. The monitor 160 processes the detector signal 254 (FIG. 2) to provide, typically, a numerical readout of the patient's oxygen saturation, a numerical readout of pulse rate, and an audible indicator or “beep” that occurs in response to each arterial pulse.\nAs shown in FIG. 2, the sensor 110 has both red and infrared LED emitters 212 and a photodiode detector 214. The monitor 160 has a sensor interface 271, a signal processor 273, a controller 275, output drivers 276, a display and audible indicator 278, and a keypad 279. The monitor 160 determines oxygen saturation by computing the differential absorption by arterial blood of the two wavelengths emitted by the sensor emitters 212, as is well-known in the art. The sensor interface 271 provides LED drive current 252 which alternately activates the red and IR LED emitters 212. The photodiode detector 214 generates a signal 254 corresponding to the red and infrared light energy attenuated from transmission through the patient finger 10 (FIG. 1). The sensor interface 271 also has input circuitry for amplification, filtering and digitization of the detector signal 254. The signal processor 273 calculates a ratio of detected red and infrared intensities, and an arterial oxygen saturation value is empirically determined based on that ratio. The controller 275 provides hardware and software interfaces for managing the display and audible indicator 278 and keypad 279. The display and audible indicator 278 shows the computed oxygen status, as described above, and provides the pulse beep as well as alarms indicating oxygen desaturation events. The keypad 279 provides a user interface for setting alarm thresholds, alarm enablement, and display options, to name a few."}
{"text": "Various methods and apparatus for performing knee arthroplasty and unicondylar knee arthroplasty in particular, are known in the art. The known methods involve resection of the tibia and femur for fitting of trial tibial and femoral implants, respectively. Once the bone has been resected and the trial implants are secured in place, the surgeon then assesses the kinematics of the knee joint. At this stage, the surgeon may transect, elevate and/or release ligaments and other soft tissue structures to achieve the desired level of deformity correction, balance in the tension of relevant ones of the ligaments and other stabilising soft tissue structures, and an acceptable range of motion of the knee joint. Additional bone resection may also be required to achieve the desired outcome. This leads to an increase in operation time with an associated increase in the risk of surgery related complications. Moreover, such additional surgical intervention following fitting of the trial implants potentially leads to subsequent increased discomfort for the patient and increased healing times. Methods and apparatus for use in arthroplasty of a knee joint are exemplified in U.S. Pat. No. 5,171,244 and U.S. Pat. No. 5,520,695."}
{"text": "When the current network multi-media broadcasting device is providing online program contents for users, the intelligent mobile terminal cannot be used to control the network multi-media broadcasting. Users are not able to use the intelligent mobile terminal to perform content menu control of the network multi-media broadcasting device. It is not possible for the users to add favorite channels and application list and set auto launch items. Furthermore, the above cannot provide users with program contents in a user-friendly manner. Typically, when users obtain programs through network TV online, a downloading bar for the buffering and loading progress will be shown, meaning that certain amount of waiting time is required to receive the information before the broadcasted program can be viewed."}
{"text": "The present invention generally relates to a device which eliminates wind rushing noise between a crash helmet and face shield structures.\nWith the ever increasing popularity of relatively high speed motorcycles, conventional protective helmets, while satisfactory to a certain degree, do not satisfy the requirements of all users of such equipment. One type of helmet commonly used is a type which protects the face of the motorcycle operator by providing a partial cylindrical transparent face shield. The periphery of the face shield overlaps the helmet shell and is fixedly secured thereto by a plurality of fasteners such as rivets or any suitable type of releasable fastener such as a screw threaded fastener, snap fastener or the like which would enable the face shield to be removed or replaced in the event of damage thereto. The bottom edge of the face shield and the crash helmet surrounding the rider's upper neck are open so as to provide access to the interior of the helmet and face shield assembly to facilitate it being placed on the head of the wearer and removed therefrom.\nWith the most common types of crash helmets equipped with transparent face shields available today, a gap or space up to one-half (1/2) inch in width exists between the periphery of the face shield where it overlaps and is affixed to the helmet with fasteners, and the underlying front leading edge of the helmet itself. As the cyclist proceeds forward, wind rushing through this gap passes down over the face and ears as it exits out of the bottom of the helmet and affixed face shield surrounding the upper neck. This wind flow disturbance behind the face shield and surrounding helmet results in eye irritation and annoying noise which increases in intensity as the cyclist goes faster and faster. The present invention eliminates this eye irritating and noisy wind flow disturbance while still allowing a milder airflow to promote comfort and help prevent fogging of the face shield."}
{"text": "Transistors such as metal oxide semiconductor field effect transistors (MOSFETs) or simply field effect transistors (FETs) are the core building blocks of the vast majority of semiconductor integrated circuits (ICs). A FET includes source and drain regions between which a current can flow through a channel under the influence of a bias applied to a gate electrode that overlies the channel. Some semiconductor ICs, such as high performance microprocessors, can include millions of FETs. For such ICs, decreasing transistor size and thus increasing transistor density has traditionally been a high priority in the semiconductor manufacturing industry. Transistor performance, however, must be maintained even as the transistor size decreases.\nA Fin field-effect transistor (FinFET) is a type of transistor that lends itself to the dual goals of reducing transistor size while maintaining transistor performance. The FinFET is a three dimensional transistor formed using a thin fin that extends upwardly from a semiconductor substrate. Transistor performance, often measured by its transconductance, is proportional to the width of the transistor channel. In a FinFET the transistor channel is formed along the vertical sidewall surfaces of the fin or on both vertical sidewall surfaces and the top horizontal plane of the fin, so a wide channel, and hence high performance, can be achieved without substantially increasing the area of the substrate surface required by the transistor.\nFinFETs provide a promising candidate for small line width technology (e.g., approximately 22 nm and below) because of their excellent short channel effect control and scalability. However, FinFETs are often formed as an array on a bulk substrate or a silicon-on-insulator (SOI) substrate in an integrated circuit, with the FinFETs densely formed on the substrates and with the substrates including a base that connects the fins. The integrate circuits that include the FinFETs often suffer from sub-fin current leakage, especially when the FinFETs are formed on a bulk substrate, whereby some of the current that is passed between a source and drain for one FinFET passes through a body of the fin to other FinFETs that are formed on the fin, thereby affecting operation of the FinFETs on the fin.\nVarious solutions to minimize or prevent current leakage have been proposed. For example, silicon-on-insulator configurations have been employed to hinder current leakage between fins. Dopant implantation has also been applied to hinder current leakage in configurations where the fins are formed in bulk semiconductor material. In particular, with dopant implantation to hinder current leakage, bulk semiconductor material between fins and at a bottom of the fins, where the fins attach to the bulk substrate, is doped with a dopant that hinders flow of current therethrough. However, dopant implantation to hinder current leakage is difficult to accurately control and may adversely impact desired operation of the FinFETs.\nAccordingly, it is desirable to provide integrated circuits that have FinFETs and methods of fabricating the integrated circuits that have FinFETs that resist sub-fin current leakage without adversely impacting desired operation of the FinFETs in the integrated circuits. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description of the invention and the appended claims, taken in conjunction with the accompanying drawings and this background of the invention."}
{"text": "1. Field of the Invention\nThis specification relates to a terminal and a wire with terminal.\n2. Description of the Related Art\nJapanese Utility Model Application No. H07-36364 discloses a terminal connected to a wire. The terminal includes a barrel to be crimped to the wire, and the barrel is formed with through holes. The terminal is crimped to the wire by fastening the barrel so that a folded part of a conductor is pressed into contact with a bottom part of the terminal by an insulation coating.\nIn the above configuration, if a cross-sectional area of the conductor becomes smaller, a force for holding the conductor is reduced and a force exerted to the insulation coating gradually increases when the wire is pulled. Then, edge parts of the through holes of the barrel easily bite into the insulation coating to tear the insulation coating. Although the barrel is crimped to the insulation coating in this way, the crimping of the insulation coating does not contribute to an improvement in the impact resistance of the wire at all and the insulation coating could not compensate for a reduction of the conductor holding force."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrophotographic photoconductor comprising an electroconductive support and a photoconductive layer formed thereon, containing lignin.\n2. Discussion of Background\nConventionally, inorganic materials such as selenium, cadmium sulfide and zinc oxide are used as photoconductive materials for electrophotographic photoconductors for use in electrophotography. Electrophotography is one of the image formation processes, through which the surface of a photoconductor is charged uniformly in the dark to a predetermined polarity, for instance, by corona charge. The uniformly charged photoconductor is exposed to a light image to selectively dissipate the electric charges of the exposed areas, so that a latent electrostatic image is formed on the photoconductor. The thus formed latent electrostatic image is developed into a visible image with a developer comprising a coloring agent such as a dye or a pigment, and a binder agent such as a polymeric material.\nFundamental characteristics required for the photoconductor for use in such electrophotography are: (1) chargeability to an appropriate potential in the dark, (2) minimum dissipation of electric charges in the dark, and (3) rapid dissipation of electric charges when exposed to light.\nHowever, while the above-mentioned inorganic materials have many advantages over other materials, they have several shortcomings from the viewpoint of practical use.\nFor instance, a selenium photoconductor, which is widely used at present, satisfies the above-mentioned requirements (1) to (3) completely, but it has the shortcomings that its manufacturing conditions are difficult and, accordingly, its production cost is high. In addition, it is difficult to work it into the form of a belt due to its poor flexibility. Furthermore, the selenium photoconductor is so vulnerable to heat and mechanical shocks that it must be handled with the utmost care.\nA cadmium sulfide photoconductor and a zinc oxide photoconductor can be respectively produced by dispersing cadmium sulfide particles and zinc oxide particles in a binder resin. However, they are poor in mechanical properties such as surface smoothness, hardness, tensile strength and wear resistance. Therefore, they cannot be used repeatedly, as they are.\nTo solve the problems of such inorganic materials, various electrophotographic photoconductors employing organic materials are recently proposed. For example, there are known a photoconductor comprising poly-N-vinylcarbazole and 2,4,7-trinitrofluorenone-9-on as described in U.S. Pat. No. 3,484,237; a photoconductor prepared by sensitizing poly-N-vinylcarbazole with a pyrylium-salt-based pigment as described in Japanese Patent Publication 48-25658; a photoconductor comprising as a main component an organic pigment as described in Japanese Laid-Open Patent Application 47-37543; a photoconductor comprising as a main component a eutectic crystal complex composed of a dye and a resin as described in Japanese Laid-Open Patent Application 47-10735; a photoconductor prepared by sensitizing a triphenylamine compound with a sensitizer pigment as described in U.S. Pat. No. 3,180,730; a photoconductor comprising an amine derivative as a charge transporting material as described in Japanese Laid-Open Patent Application 57-195254; a photoconductor comprising poly-N-vinylcarbazole and an amine derivative as charge transporting materials as described in Japanese Laid-Open Patent Application 58-1155; and a photoconductor comprising as a photoconductive material a polyfunctional tertiary amine compound, in particular, a benzidine compound, as described in U.S. Pat. No. 3,265,496, Japanese Patent Publication 39-11546 and Japanese Laid-Open Patent Application 53-27033.\nThese organic electrophotographic photoconductors are still unsatisfactory for use in practice, especially with respect to the durability thereof. As the demand for a photoconductor with higher durability is increasing year by year, ensuring the charging stability in the photoconductor has become a requirement that cannot be ignored. When a photoconductor with a decreased charging stability is employed, in the case of an electrophotographic copying machine, the lowering of image density is caused; and in the case of a laser printer employing the reversal development method, image quality is lowered, for instance, with the occurrence of toner deposition on the background of printed images. In order to solve these problems, it is proposed to provide an intermediate layer between an electroconductive support and a photoconductive layer. However, in the case where a material with high resistivity, having high barrier properties, is employed for the intermediate layer, the photosensitivity of the photoconductor decreases and the residual potential thereof increases although the charging characteristics are improved. When a material with relatively low resistivity is employed for the intermediate layer to prevent the increase of the residual potential, the charging stability of the photoconductor is still insufficient.\nWhen a photoconductor is incorporated into a copying machine in practice, the photoconductor is exposed to ozone generated from a corona charger. By the ozone thus generated, organic materials such as a charge transporting material contained in a photoconductive layer of the photoconductor are oxidized, so that there are caused the problems that the photosensitivity is decreased, the residual potential is increased, and the charging potential is decreased.\nTo solve such problems, the addition of an antioxidant to the photoconductive layer is proposed as in Japanese Laid-Open Patent Applications 57-122444, 50-33857, 63-18355, 63-18356 and 3-172852; and the provision of a gas-barrier resin layer on the charge transport layer is also proposed as in Japanese Laid-Open Patent Application 63-135955.\nThe above-mentioned countermeasures, however, have not yet solved the problems of the conventional photoconductors, such as the increase of the residual potential, the lowering of the photosensitivity, and the insufficient improvement of the durability thereof."}
{"text": "1. Field of the Invention\nThis invention is directed to the field of large data input/output (I/O) transfer into, and out of, large compute systems and large data storage systems. It is particularly directed towards reducing the data access latency times and increasing the data I/O transfer rates for compute clients processing extremely large data sets within large, shared-memory compute systems.\n2. Description of Prior Art\nDemand is growing to find solutions to the “big data” challenge faced by government and business, for example, in areas such as fraud detection, remote sensor data analysis, financial modeling, social network analysis, and risk analysis. Many data intensive computing applications need to process an enormous amount of data in a timely manner, and if too much time is spent loading a required data set into a compute client's memory, then the entire production workflow will be delayed. The same logic holds for processes that are generating very large data sets. The newly generated large data sets need to be quickly made available for the next step in the larger production workflow. To support such an environment, a supercomputer architecture is required that incorporates hundreds or thousands of CPU cores for handling multiple processing threads and supersized memories to hold more of the very large data sets. Simply put, the more data the supercomputer can process at any given time, the quicker the results are presented to decision makers. When it comes to mission-critical pattern discovery problems in government and business, whoever obtains the answers first is one step ahead of their adversary or competitor.\nThe challenges to meeting “big data” demands are many:\nThe compute hardware needs to be able to scale up to an extraordinary level—thousands of CPU cores and multiple terabytes of memory—which is far beyond the capacity of any commodity X86 server. Computer software needs to be able to take advantage of the scaled-up hardware platform. And, an extremely efficient data access input/output system that can quickly load the “big data” into system memory and just as quickly store the processed data onto reliable media is needed to complete the system.\nThe ability to build larger and larger compute systems has been facilitated by advances in massively parallel compute system techniques, algorithmic resource orchestration techniques (Map Reduce and others), low-cost/high-performance CPU's, and the advent of both open-source and proprietary Parallel File Systems. Parallel File Systems such as “Lustre” [See “Lustre scalable storage” Copyright 2006 Cluster File Systems —Rights owned by ORACLE Corp.], and IBM's GPFS (General Parallel File System) [Refer to F. B. Schmuck and R. L Haskin, GPFS: a shared-disk file system for large computing clusters, in proceedings of Conference of Files and Storage Technologies (FAST'02), 2002] have allowed these newer and larger parallel compute systems to scale up to higher and higher numbers of compute nodes and compute-cores and still have enough File System I/O Bandwidth to allow for very good compute performance on many classes of applications.\nParallel File Systems generate very high “Total I/O” Bandwidth levels by utilizing multiple “File System Data Storage Servers” like the dual socket server 110 in FIG. 1 (Lustre uses “Object Storage Servers (OSS's) and IBM's GPFS uses NSD Servers (Network Shared Disk Servers)). Each File System Data Storage Server (OSS or NSD Server) is attached to one or more “Data Storage Devices,” like 109 in FIG. 1, via a storage network fabric 106 that in many HPC environments is implemented using Infiniband (IB) switches and (IB) Host Channel Adapters (HCA's) 104. The “Data Storage Device” 109 often consists of a “Storage Controller” 107, and a High-Performance or High-Capacity Disk Storage Devices 108 or Solid-State Disk (SSD) Storage Devices (many new systems are using SSDs in combination with SATA or SAS Disks to balance total I/O operations with total storage capacity). Sometimes “File System Data Storage Servers,” like 110 in FIG. 1, are paired with Storage Subsystems, like 109 in FIG. 1, that are implemented as raw “Trays-of-Disks” without Disk-Controllers. These Trays-of-Disks are referred to as JBOD's (just a Bunch Of Disks), and require additional Data-Block management software to run within the “File System Data Storage Server” 110 to accomplish the functionality required for the “JBOD” to serve as a suitable Storage Device 109. The “Zetta File System,” or ZFS [open source licensed By Oracle Corp.], is one candidate application for handling JBOD's and is being utilized by the Lawrence Livermore National Labs (LLNL) team, in combination with the Lustre File System as a part of their “Next-Generation” Sequoia Super Computer installation which hopes to have a “Total File System I/O Bandwidth of between 500 GB/s and 1000 GB/s.\nThe data blocks in the Parallel File System are “Striped” or spread across multiple “File System Data Storage Servers,” like the 111 grouping of FIG. 1, and their associated collective sets of Storage Devices 113 in FIG. 1. The sequence of events for a compute client to receive the data stored in the Parallel File Systems after a client data request is:\n1—The data blocks are copied from “Disk Storage Devices —108 data storage locations” and into “Storage Controllers 107 memory—Data Movement #1,\n2—The data blocks are then copied from “Storage Controller 107 memory” into “File System Data Storage Servers 110 memory”—Data Movement #2,\n3—The data blocks are then copied from all of the “File System Data Storage Servers 110 memory” into the client 101-1 memory space—Data Movement #3.\nMoving the data blocks three times before they can be utilized is very inefficient, but for many parallel applications the inefficiency is accepted as current practice and the job schedules are allocated by compute time required as well as data I/O transfer time. All of these “Data Movement processes” are running in parallel, simultaneously across all of the File System Data Storage Servers and their respective Storage Controllers and Disk Storage Devices. The aggregate data I/O speed for the entire File System is a function of the available bandwidth for all of the File System Data Storage Servers operating in unison. The grouping of File System Data Storage Servers 111 in FIG. 1 has the potential I/O bandwidth of 8 GB/s because each of the 4 servers can maintain 2 GB/s and the set of 2 disk Storage Devices, 109 in FIG. 1, can each maintain 4 GB/s.\nThe total I/O performance of the Parallel File Systems can be scaled up by adding more File System Data Storage Servers, see 201 in FIG. 2, and more Disk Storage Devices, 202 in FIG. 2. The additional server and storage resources allow the example Parallel File System in FIG. 2 to provide an aggregate potential data I/O bandwidth of 20 GB/s instead of the original potential of I/O speed of 8 GB/s. The combination of adding more Storage Servers and more Storage Devices creates a very high Total potential I/O Bandwidth solution for use by the massively parallel compute systems and their applications.\nThe use of Parallel File Systems, like Lustre and GPFS, for massively parallel compute systems has worked well because each compute node in the massively parallel compute system typically needs to use only a small portion of the I/O Bandwidth to accomplish its compute tasks and to deliver intermediate or final results. The High Speed Parallel File Systems in use at Large Supercomputing Centers are used to take very rapid “Snap-shots” of intermediate compute results for complex compute jobs that can sometimes take days or weeks to complete. The snapshots are called “Check Points” and they require very High Total I/O Bandwidth since each of the many thousands of compute nodes are sending copies of their intermediate compute results out for analysis and algorithmic tuning.\nOne of the HPC systems at Oak Ridge National Laboratory, Jaguar, has over 26,000 compute nodes with approximately 180,000 compute-cores, and its Lustre Parallel File Systems operate at 240 GB/s and 44 GB/s. That meant that each “File System Client”, like 101-2 in FIG. 1, on a small compute node, could “Simultaneously Share” about 9-12 MB/s on average of the “Total File System I/O,” when all of the many thousands of compute nodes were working together in parallel on one large problem. The peak File System I/O performance for a single specific “File System Client,” like 101-2 in FIG. 1, on a Jaguar compute node was just over 1.25 GB/s. This means that the peak “File System Client” I/O Bandwidth performance available for a single node on Jaguar is 0.5% of the “Total I/O” Bandwidth available from the entire File System. This level of I/O performance works very well for the large scale scientific simulation problems that have been tailored to work on the massively scaled parallel compute systems like “jaguar” with it's thousands of separate compute nodes that typically begin processing jobs with a few small-sized data sets for each compute node. But there are many additional types of compute problems that require the manipulation and processing of very large data sets that must first be loaded into the memory space of one compute system before processing can begin.\nMany US Government projects have requirements to rapidly access and process large numbers of files in the 5 GB to 15 GB range. There are many fielded systems are producing data sets at a rate of 1 Tera Byte (TB) or more per hour, and there are several ongoing projects that produce data sets that range from 4 to 6 TB in size. The users of these large, Multi-TB data sets would like to have the ability to rapidly load and process entire multi-TB files in “Large Shared Memory” compute systems, like 112 in FIGS. 1 and 2, and utilize the compute systems 100's or 1000's of compute-cores to reduce the “Raw multi-TB” data sets into useful, user friendly, result sets. The task of crafting algorithms to process these Very Large Multi-TB data sets is much easier to accomplish if the entire data set can fit within the internal memory address space of a large “Common Global Shared Memory” compute system and be acted on by all of the 100's or 1000's of compute-cores resident with the single large “Common Global Shared Memory” compute system.\nAn example of using a “Common Global Shared Memory” compute system in combination with a typical Parallel File System to process very large data sets can be further examined by using the performance characteristics of Oak Ridge's “Jaguar” Supercomputer. Jaguar's “Peak I/O Performance” for a single compute node was 1.25 GB/s. Using that I/O Bandwidth value for the I/O performance of a single “File System Client” within a “Large Shared Memory” compute system, such as 112 in FIG. 1 or 2, results in a 15 GB file loading into the “Large Shared Memory” Compute system in just 12 seconds. An acceptable data I/O transfer time for many situations. But the time required to load a 6 TB file from a modular storage device or a current implementation of a typical Parallel File System, 111 and 201 and 202 from FIG. 2, into the same large “Common Global Shared Memory” compute system, 101-1 in 112 of FIG. 2, would be at least 4,800 seconds! Having 100's or 1000's of processing cores waiting over 4,000 seconds for a data set to load into a system like 112 in FIG. 1 or 2, would be a terrible waste of an expensive compute resource.\nThe “Peak I/O bandwidth” available for a specific application compute client from the total aggregate I/O provided by the Parallel File Systems' entire collection of “File System's Data Storage Servers,” like 111 and 201 added together in FIG. 2, can be limited by many of the physical and logical elements that interconnect a typical Lustre or GPFS Parallel File System implementation.\nA significant implementation element that can limit the “I/O Bandwidth Performance” for a specific “File Systems' Client” is the total number, and bandwidth capacity, of the physical or logical I/O network pathways available to link the specific “Application Compute Client” with the “Data Storage Servers” of the Parallel File System. The designs and implementations of “Application Compute Clients” will vary across the range of commercially available Parallel File Systems. The number of I/O network pathways supported by a specific “Application compute client” implementation, combined with the I/O bandwidth available for “Application compute client” use within each supported I/O network pathway, will significantly affect the total I/O Bandwidth for the specific “File System Client.”\nThe physical linkage between the “Application compute client” processes 101-1 and the collective set of “File System Data Storage Servers” 111 and 201, is depicted in FIGS. 1 and 2 with the HCA's 104 supporting the external I/O connectivity from the 101-1 Client within the large “Common Global Memory Compute System,” 112, and the Infiniband Switch 105 that in-turn provides the I/O network pathway connectivity to the entire collection of “File System Data Storage Servers” 111 and 201 in FIG. 2.\nOne specific current example of how a Parallel File System's design and implementation of its “Application compute client” can limit the total amount of File System I/O Bandwidth available to a specific “Application compute client” is Lustre's current Client design and implementation. As currently implemented, a Lustre Client can have one physical Infiniband connection. This means that a QDR (Quad Data Rate=4 GB/s) or DDR (Double Data Rate=2 GB/s) IB connection between a Lustre File System Client 101-1 and HCA 104 in FIG. 2 over to the 105 IB switch, would be the only I/O pathway between the Lustre File System Client and the collection of “File System Data Storage Servers” (OSS's for Lustre). Current Lustre Client instances utilizing a single QDR IB connection as its I/O network pathway to the “File System Data Storage Servers” have been able to achieve sustained I/O rates of up to 3.2 GB/s after considerable tuning adjustments were made, and while the Lustre Client was servicing several separate processes within a “Large Shared Memory” compute system with data blocks from several unique files that were striped across all of the “File System Data Storage Servers” (OSS servers) within the specific Lustre File System.\nA redesign and re-implementation of the Lustre Client software would be required to permit the Lustre “File System Client” process to utilize two or more physical IB connections as its I/O network pathway to the “File System Data Storage Servers.” An implementation that allowed a Lustre File System Client to utilize two QDR IB connections would potentially be able to achieve I/O bandwidth rates up to a maximum of 6 or 7 GB/s, but even this level of File System Client I/O Bandwidth performance is still a very small percentage of the total I/O Bandwidth available from Parallel File Systems that can achieve 100's of GB/s of total I/O Bandwidth. If the File System Client 101-1 in 112 of FIG. 2 was able to operate at 6 GB/s of I/O Bandwidth, it would still be utilizing less than ⅓ of the total potential File System I/O of 20 GB/s derived from the Parallel File System's “File System Data Storage Servers” in 111 and 201 in FIG. 2.\nExisting Parallel File Systems (like Lustre and GPFS) perform well at distributing a small share of the total file system I/O to the thousands of compute nodes in today's large supercomputing compute clusters. However, for cases where very large data files need to be quickly loaded into a single, multi-processor compute node, the existing Parallel File Systems are not capable of providing more than a very small percentage of the total I/O Bandwidth to an individual compute client. A more efficient use of the Large Shared Memory Compute System, 112 in FIG. 2, would be had if there were a way for all of the Parallel File Systems aggregate I/O to be available for the compute client or clients within the Large Shared Memory Compute system.\nTherefore, there is a need for a solution to improve the “Peak I/O Bandwidth Performance” for a single “File System Client,” 101-1 in FIGS. 1 and 2, in a large “Common Global Shared Memory” Compute System, 112 in FIGS. 1 and 2."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of soybean breeding. In particular, the invention relates to the novel soybean variety 01056961.\n2. Description of Related Art\nThere are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include higher seed yield, resistance to diseases and insects, better stems and roots, tolerance to drought and heat, better agronomic quality, resistance to herbicides, and improvements in compositional traits.\nSoybean, Glycine max (L.), is a valuable field crop. Thus, a continuing goal of plant breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties.\nThe oil extracted from soybeans is widely used in food products, such as margarine, cooking oil, and salad dressings. Soybean oil is composed of saturated, monounsaturated, and polyunsaturated fatty acids, with a typical composition of 11% palmitic, 4% stearic, 25% oleic, 50% linoleic, and 9% linolenic fatty acid content (“Economic Implications of Modified Soybean Traits Summary Report,” Iowa Soybean Promotion Board & American Soybean Association Special Report 92S, May 1990)."}
{"text": "A known semiconductor device of this kind is depicted and described, for example, in patent DE 195 06 093 02. A schematic diagram of such a known semiconductor device is provided in FIG. 2. In that case, a microcooler 20 is fabricated by bonding together plural copper foils structured by etching. The individual layers form in conjunction a coolant inlet 24; a cooling channel 26, which conducts the coolant to the region of the microcooler 20 on which a power laser bar 12 is mounted; and a coolant outlet 28. The coolant flows along the arrow 30 from the inlet 24 to the outlet 28. Implemented in at least one region 32 are microstructures: narrow channels, for example. Particularly effective heat exchange takes place in that region because of turbulent flow of the coolant.\nThe laser bar 12 is soldered to the front edge of the microcooler by means of a soft solder 52—indium, for example. Mounting the bar 12 directly on the copper block 20 improves heat transfer from the laser bar to the cooler."}
{"text": "1. Field of the Invention\nThis invention relates to a portable, hand-held, device for transferring data to and from a data processing system via a graphical user interface, and, more particularly, to a hand-held data storage device that generates graphical user input signals to facilitate data transfer.\n2. Description of the Prior Art\nIn the prior art, there are a variety of devices and networks to move data from one data processing system (e.g., a personal computer system) to another data processing system. These include diskettes (e.g., magnetic and optical), local area hardwired networks, various wireless transmission networks, and semi-conductor memory cards.\nMore specifically to this invention, U.S. Pat. No. 4,102,493 issued to Moreno et al. (Moreno '493) describes a method and apparatus for transporting data from one device to another. In Moreno '493, data is transferred from a processor system to a portable card in which data is electronically stored, modified, and transferred to another processor system. Moreno '493 transports data by means of a hand-held device, but the device does not provide a user interface to facilitate the transfer (upload or download) of the data.\nU.S. Pat. No. 4,125,871, issued Nov. 14, 1987, to Martin et al. (Martin '871) describes a portable data entry device in which a user keys in data, which is stored in the device until a later time when the data is uploaded to a computer system. The portable data entry device is wholly contained within a small housing. The device includes an electronic memory capable of storing a plurality of multiple character records and includes manually operable controls for sequencing through the memory for review and updating of previously entered data. A connector is provided on the housing by which the device can be directly connected to a data system for the readout of the stored data. The device is self powered and contains circuitry operative to conserve available energizing power.\nU.S. Pat. No. 4,689,757, issued to Downing et al. (Downing '757), describes an apparatus for transporting information captured at a coin counter to a computer. The system is comprised of a discrete machine event counting module which records and stores a count of machine operation and can include means for recording the time of some selected event or events. The module also stores an identification code for the particular machine. The module can be connected directly through a microprocessor to a central processing center, or it can be located at a machine or at a group of machines.\nThe transfer unit can then be transported to access means for the central processing center and the information that was obtained from the module will be transferred to centers for processing and tabulation.\nOf course, graphical user interfaces which allow a user to generate computer commands by means of a graphical icon display and display pointer are well known and widely used in the art (see, for example, U.S. Pat. No. 5,204,947 and the materials referenced therein). Notwithstanding, in the prior art there is no generally applicable, simple way to transfer data among diverse systems."}
{"text": "1. Field of the Invention\nThe present invention relates to light-emitting modules for light communication, and in particular to a light-emitting module which is preferably used as a pig tail type in which an optical fiber is fixed.\n2. Description of the Related Art\nIn a pig tail type of light-emitting module in which a laser diode, a lens and an optical fiber are incorporated, a laser beam emitted from the laser diode is converged by the lens and the laser beam converged by the lens and then optically coupled with the optical fiber. Such optical coupling requires not only essential accurate alignment in directions (X- and Y-directions) perpendicular to the optical direction of the laser beam in the optical fiber, but also accurate alignment in the direction (Z-direction) of the optical axis due to variations in the position of the laser diode chip when the chip is manufactured. Therefore, conventionally, the laser diode and the lens are aligned at first so as to be incorporated, and subsequently, the three components being the laser diode, the lens and the optical fiber are incorporated by adjusting the optical fiber in the X-, Y- and Z-directions.\nHowever, according to the above conventional light-emitting module, the optical fiber must be tri-axially adjusted in the X-, Y- and Z-directions, which requires not only a complicated adjusting mechanism but also considerable time for adjustment. This greatly increases manufacturing cost. In general, a laser diode with a cap has various chip position shifts of approximately .+-.60 .mu.m in the X-, Y- and Z-directions, thus, if adjustment in the Z-direction is omitted, coupling efficiency greatly deteriorates, which disadvantageously causes deterioration in performance of the light-emitting module."}
{"text": "1. Field of the Invention\nThe present invention concerns the fields of cancer diagnostics and targeted delivery of therapeutic agents to cancer cells. More specifically, the present invention relates to compositions and methods for identification and use of peptides that selectively target cancer cell receptors, such as the IL-11 receptor and/or the GRP78 receptor. In particular embodiments, the targeted receptors are preferentially expressed in prostate cancer, especially in metastatic prostate cancer. In certain embodiments, the invention concerns compositions and methods of use of novel phage-based gene delivery vectors.\n2. Description of Related Art\nTherapeutic treatment of many disease states is limited by the systemic toxicity of the therapeutic agents used. Cancer therapeutic agents in particular exhibit a very low therapeutic index, with rapidly growing normal tissues such as skin and bone marrow affected at concentrations of agent that are not much higher than the concentrations used to kill tumor cells. Treatment of cancer and other organ, tissue or cell type confined disease states would be greatly facilitated by the development of compositions and methods for targeted delivery to a desired organ, tissue or cell type of a therapeutic agent.\nRecently, an in vivo selection system was developed using phage display libraries to identify targeting peptides for various organs, tissues or cell types in a mouse model system. Phage display libraries expressing transgenic peptides on the surface of bacteriophage were initially developed to map epitope binding sites of immunoglobulins (Smith and Scott, 1986, 1993). Such libraries can be generated by inserting random oligonucleotides into cDNAs encoding a phage surface protein, generating collections of phage particles displaying unique peptides in as many as 109 permutations. (Pasqualini and Ruoslahti, 1996, Arap et al, 1998a; Arap et al 1998b).\nIntravenous administration of phage display libraries to mice was followed by the recovery of phage from individual organs (Pasqualini and Ruoslahti, 1996). Phage were recovered that were capable of selective homing to the vascular beds of different mouse organs, tissues or cell types, based on the specific targeting peptide sequences expressed on the outer surface of the phage (Pasqualini and Ruoslahti, 1996). A variety of organ and tumor-homing peptides have been identified by this method (Rajotte et al., 1998, 1999; Koivunen et al., 1999a; Burg et al., 1999; Pasqualini, 1999). Each of those targeting peptides bound to different receptors that were selectively expressed on the vasculature of the mouse target tissue (Pasqualini, 1999; Pasqualini et al., 2000; Folkman, 1995; Folkman 1997). Tumor-homing peptides bound to receptors that were upregulated in the tumor angiogenic vasculature of mice (Brooks et al., 1994b; Pasqualini et al., 2000). In addition to identifying individual targeting peptides selective for an organ, tissue or cell type (Pasqualini and Ruoslahti, 1996; Arap et al, 1998a; Koivunen et al., 1999b), this system has been used to identify endothelial cell surface markers that are expressed in mice in vivo (Rajotte and Ruoslahti, 1999).\nThis relative success notwithstanding, cell surface selection of phage libraries has been plagued by technical difficulties. A high number of non-binder and non-specific binder clones are recovered using previous in vivo methods, particularly with components of the reticuloendothelial system such as spleen and liver. Removal of this background phage binding by repeated washes is both labor-intensive and inefficient. Cells and potential ligands are frequently lost during the many washing steps required. Methods that have been successful with animal model systems are unsatisfactory for identifying human targeting peptides, which may differ from those obtained in mouse or other animal model systems.\nAttachment of therapeutic agents to targeting peptides has resulted in the selective delivery of the agent to a desired organ, tissue or cell type in the mouse model system. Targeted delivery of chemotherapeutic agents and proapoptotic peptides to receptors located in tumor angiogenic vasculature resulted in a marked increase in therapeutic efficacy and a decrease in systemic toxicity in tumor-bearing mouse models (Arap et al., 1998a, 1998b; Ellerby et al., 1999). However, the targeted delivery of anti-cancer agents in humans has not yet been demonstrated. The targeted receptors reported in previous studies may be present in angiogenic normal tissues as well as in tumor tissues and may or may not be of use in distinguishing between normal tissues, non-metastatic cancers and metastatic cancer. A need exists for tumor targeting peptides that are selective against human cancers, as well as for targeting peptides that can distinguish between metastatic and non-metastatic human cancers.\nAttempts have been made to target delivery of gene therapy vectors to specific organs, tissues or cell types in vivo. Directing such vectors to the site of interest would enhance therapeutic effects and diminish adverse systemic immunologic responses. Adenovirus type 5 (Ad5)-based vectors have been commonly used for gene transfer studies (Weitzman et al., 1997; Zhang, 1999). The attachment of Ad5 to the target cell is mediated by the capsid's fiber knob region, which interacts with cell surface receptors, including the coxsackie adenovirus receptor (CAR) and possibly with MHC class I (Bergelson et al., 1997; Hong et al., 1997). Upon systemic administration in vivo, binding of virus to CAR can result in unintended enrichment of vectors in non-targeted but CAR-expressing tissues. Conversely, target cells that express little or no CAR are inefficiently transduced. A need exists to develop novel gene therapy vectors to allow more selective delivery of gene therapy agents.\nA need also exists to identify receptor-ligand pairs in organs, tissues or cell types. Previous attempts to identify targeted receptors and ligands binding to receptors have largely targeted a single ligand at a time for investigation. Identification of previously unknown receptors and previously uncharacterized ligands has been a very slow and laborious process. Such novel receptors and ligands may provide the basis for new therapies for a variety of disease states, such as cancer and/or metastatic prostate cancer."}
{"text": "The present invention relates to a pressure jet cleaning appliance with a pressure pump driven by a driving means from whose pressure pipe a bypass pipe, which, when the pressure pipe is closed, is opened by means of a closing valve provided with a closing element, leads to the suction pipe of the pressure pump and having a control element that is connected to a governor of the driving means.\nIn non-stationary operating pressure jet cleaning appliances, an electromotor or an internal combustion engine is frequently employed as driving means of the pressure pump that is operated at a constant rotational speed. In the case of the electromotor, the rotational speed control is effected independently of the load by means of the fixed number of poles and the constant mains frequence and, in the case of the internal combustion engine, by means of a rotational speed-dependent throttle valve position. In such a case the control of the system is effected in that, with the opened spray gun, the driving means drives the pressure pump with the power that is necessary for the operation of the pressure jet cleaning appliance. When the spray gun is then closed with the aid of the normally provided control handle, a bypass closing valve enters into action in such a way that, when the pressure pipe is closed by the spray gun, a bypass pipe to the suction pipe opens so that the pressure pump then delivers the cleaning fluid into the circulation. The power of the pressure pump that is necessary in this bypass operation is considerably less than the power required in the normal spraying, squirting or cleaning operation.\nIn order to reduce the noise development and the wear during the operating intervals (bypass operation), a control is expediently effected in that, when the pressure pipe is closed, the driving means is regulated down to a rotational speed that is lower than the nominal operating speed.\nFor this, in a known high-pressure jet cleaning appliance, a pressure sensor is provided in the pressure pipe which detects the pressures which occur inside the pressure pipe, in particular the lowering of the pressure which occurs after the closing of the pressure pipe and the opening of the bypass valve connected therewith, and the increase in the fluid pressure which occurs in the pressure pipe following the reopening of the pressure pipe and reduces the rotational speed of an internal combustion engine provided as driving means as soon as the pressure in the pressure pipe, due to the opening of the bypass valve, drops below a certain value and, when a certain pressure value is exceeded, once more adjusts the throttle control rod of the internal combustion engine to an increase in the rotational speed (DE 39 02 252 C1).\nHowever, this method of regulation is subject to various disadvantages. For it has to be stated that such a control system operates according to quite specific pressure changes which arise in the pressure pipe under quite specific directions subsequent to the opening or closing of the bypass valve, said changes being then utilized for the control. This means that the control is merely able to lag timewise behind the pressure conditions occurring, that is to say, the control becomes only effective when a rise in pressure or a drop in the pressure has taken place.\nFurthermore, in pressure jet cleaning appliances there arises the problem that an operation of the pressure jet cleaner is only possible at one pressure, viz. the maximum pressure, independently of the circumstance whether this pressure is actually required. That is why it is desirable to adapt the pressure of the squirted-out cleaning medium as well as the driving power needed herefor to the case of application. With the DE 39 22 956 A1, a high-pressure jet cleaning device has already been proposed which comprises a high-pressure pump, a pressure pipe leading from the latter to a closable outlet, a bypass pipe leading from the same to the induction side of the pump and a closing valve that closes the bypass pipe. The device makes an automatic pressure discontinuation of the pump possible when closing the outlet, is constructed in such a way that, upstream of the outlet and downstream of the branching of the bypass pipe, a flow controller is provided which closes the closing valve of the bypass pipe when the outlet is closed. In this case the closing valve is constructed in the form of a closing valve and provision has been made for an adjustable stop, with the aid of which the closing motion of the closing valve can be limited. However, with such a set-up, only the volume flow is regulated.\nThis invention is based upon the technical problem of constructing, in a pressure jet cleaning appliance of the type stated in the beginning, a control in such a way that the same is independent of pressure changes that occur due to the closing or opening of the closing valve and, more particularly, prior to the occurrence of these pressure changes, adjusts the rotational speed of the driving means."}
{"text": "1. Field of the Invention\nThe present invention relates to a level gauge for detecting a level of liquid helium which is accommodated in a container made of metals, glasses or other materials. More particularly, the invention relates to a level gauge for detecting a level of liquid helium which uses, as a sensing element, a wire made of an amorphous superconductive alloy obtained by rapid by quenching a molten alloy. The level of liquid helium is detected by measuring an electric current, voltage and/or electric resistance of the sensor element.\n2. Description of Related Technology\nLiquid helium level gauges that use superconducting alloy sensing elements rely on the electrical resistance changes of the element to indicate liquid level. The portion of the element submerged in the liquid becomes superconductive, i.e. no resistance to electrical current. The portion above the liquid is not superconductive and resists electrical flow at a constant rate over its length. If the sensing element is homogeneous, has a constant width, and has a constant thickness, the resistance properties will be constant over the length of the sensor element. By passing an electrical current through the submerged element, measuring the electrical current, and comparing the value to a calibration relationship, the level of the helium can be determined.\nIn the level gauge disclosed in U.S. Pat. No. 4,655,079 (which is herein incorporated by reference), the superconductive alloy is represented by the following formula: EQU Zr.sub.100-x (Ru.sub.y Rh.sub.1-y)x\nwherein x represents the contents of Ru and/or Rh in aomic % and has a numeral value of 22.5<x<27.5; and y represents a numerical value of 0<y<1.\nHowever, the superconducting transition temperature (Tc) of that superconductive alloy ranges from 4.2K to 4.5K. This transition temperature is quite close to the temperature of liquid helium (4.2K). As the pressure in the storage vessel changes, the accuracy of the level measurements can decrease."}
{"text": "A known semiconductor light emitting apparatus can include a substrate, a frame body mounted on the substrate, the frame body having a light-shielding property and a reflective property, a semiconductor light emitting element mounted on the substrate that is exposed within the frame body, and a light-transmitting resin filled inside the frame body to resin-seal the semiconductor light emitting element. This type of conventional semiconductor light emitting apparatus can be manufactured by the method shown in FIGS. 1A to 1E, for example, including:\nStep (a): fabricating a multi-piece substrate 101 on which a plurality of semiconductor light emitting elements 100 are mounted at predetermined intervals (FIG. 1A);\nStep (b): performing primary transfer molding using a light transmitting resin to form a light transmitting portion 102 with a predetermined thickness on the mounted substrate 101 (thereby resin-sealing semiconductor light emitting elements shown in FIG. 1B);\nStep (c): cutting the light transmitting portion 102 using a dicer to form grooves 13 at predetermined intervals in the longitudinal and transverse directions with a predetermined depth and a predetermined width (FIG. 1C);\nStep (d): performing secondary transfer molding using a resin having a light-shielding property and a reflective property to form frame bodies 104 within the grooves (FIG. 1D);\nStep (e): cutting the product along the frame bodies 104 using a dicer to separate individual substrates 101 having the respective semiconductor light emitting elements mounted thereon together with the respective frame bodies 104 (FIG. 1E).\nThrough the above steps (a) through (e), a plurality of semiconductor light emitting elements 105 having the appearance shown in FIG. 1F can be simultaneously completed (for example, Japanese Patent Application Laid-Open Nos. 2002-344030 and 2003-31854).\nAnother method has been proposed as follows (for example, see Japanese Patent Application Laid-Open No. 2006-324623). First, a frame body is formed by performing primary transfer molding. Then, the product is released from the mold, and a light emitting body is disposed within the frame body. The resulting product is again set in the mold to perform secondary transfer molding so that a space within the frame body is filled with a light-transmitting resin for resin-sealing."}
{"text": "The steel for machine structural use such as a gear, shaft, pulley, constant velocity joint and the like utilized for a variety of gear transmission devices, to begin with a transmission for an automobile and a differential gear, as well as a crank shaft, con'rod and the like, is generally finished into a final shape by performing the work of forging and the like and thereafter performing cutting work. Because the cost required for the cutting work occupies a major portion in the manufacturing cost, steel material constituting the steel for machine structural use is required to be excellent in machinability. Therefore, technologies for improving machinability have been disclosed from the past.\nThe typical examples of such technologies are to add Pb and to form MnS by adding S. However, because Pb is hazardous for a human body, its use has come to be restricted. Also, with respect to the parts in which deterioration of the mechanical property caused by sulfide becomes a problem, there is a limit in using S. Further, in cutting work of a gear and the like particularly, gear cutting with a hob is generally performed, however, the cutting in this case differs from the continuous cutting such as what is called lathe turning but is in a manner called as the intermittent cutting. At present, steel material improving the machinability in hobbing has been scarcely materialized. The tool raw material used for a hob is a high-speed steel, and is generally performed with coating of TiAlN and the like. In this case, it is known that the tool surface wears while being oxidized by repeating cutting and idle rotation in working under a comparative low speed.\nAs a method for improving the intermittent cutting machinability, in the patent document 1, steel material is described which is excellent in the intermittent cutting machinability (tool life) under a high speed (cutting speed: 200 m/min or more) by containing Al: 0.04-0.20% and O: 0.0030% or less.\nIn the patent document 2, a steel for machine structural use is described which contains C: 0.05-1.2%, Si: 0.03-2%, Mn: 0.2-1.8%, P: 0.03% or less, S: 0.03% or less, Cr: 0.1-3%, Al: 0.06-0.5%, N: 0.004-0.025%, O: 0.003% or less respectively, contains Ca: 0.0005-0.02% and Mg: 0.0001-0.005% with solid solution N: 0.002% or more in steel, the remainder being iron and unavoidable impurities, and satisfies\n(0.1×[Cr]+[Al])/[O]≦150.\nAlso, in the patent document 3, a steel for machine structural use is described which contains C: 0.1-0.85%, Si: 0.01-1.0%, Mn: 0.05-2.0%, P: 0.005-0.2%, total Al: exceeding 0.1% and 0.3% or less, total N: 0.0035-0.020%, with solid solution N being limited to 0.0020% or less."}
{"text": "Adiponitrile (ADN) is a commercially important and versatile intermediate in the industrial production of nylon polyamides useful in forming films, fibers, and molded articles. ADN may be produced by hydrocyanation of 1,3-butadiene (BD) in the presence of transition metal complexes comprising various phosphorus-containing ligands. For example, catalysts comprising zero-valent nickel and monodentate phosphorus-containing ligands are well documented in the prior art; see, for example, U.S. Pat. Nos. 3,496,215; 3,631,191; 3,655,723 and 3,766,237; and Tolman, C. A., McKinney, R. J., Seidel, W. C., Druliner, J. D., and Stevens, W. R., Advances in Catalysis, 1985, Vol. 33, pages 1-46. Improvements in the hydrocyanation of ethylenically unsaturated compounds with catalysts comprising zero-valent nickel and certain multidentate phosphite ligands are also disclosed; e.g., see: U.S. Pat. Nos. 5,512,696; 5,821,378; 5,959,135; 5,981,772; 6,020,516; 6,127,567; and 6,812,352.\n3-Pentenenitrile (3PN) may be formed through a series of reactions as illustrated below.\n\nAccording to abbreviations used herein, BD is 1,3-butadiene, HC≡N is hydrogen cyanide, and 2M3BN is 2-methyl-3-butenenitrile. A method to increase the chemical yield of 3PN from BD hydrocyanation includes the catalytic isomerization of 2M3BN to 3PN (Equation 2 above) in the presence of NiL4 complexes as disclosed in U.S. Pat. No. 3,536,748. Co-products of BD hydrocyanation and 2M3BN isomerization may include 4-pentenenitrile (4PN), 2-pentenenitrile (2PN), 2-methyl-2-butenenitrile (2M2BN), and 2-methylglutaronitrile (MGN).\nIn the presence of transition metal complexes comprising various phosphorus-containing ligands, dinitriles such as ADN, MGN, and ethylsuccinonitrile (ESN) may be formed by the hydrocyanation of 3PN and 2M3BN, as illustrated in Equations 3 and 4 below. Equation 4 also shows that 2M2BN can be formed when 2M3BN undesirably isomerizes in the presence of a Lewis acid promoter that may be carried over from a pentenenitrile hydrocyanation reaction zone.\n\nThe hydrocyanation of activated olefins such as conjugated olefins (e.g., 1,3-butadiene) can proceed at useful rates without the use of a Lewis acid promoter. However, the hydrocyanation of un-activated olefins, such as 3PN, require at least one Lewis acid promoter to obtain industrially useful rates and yields for the production of linear nitriles, such as ADN. For example, U.S. Pat. Nos. 3,496,217, 4,874,884, and 5,688,986 disclose the use of Lewis acid promoters for the hydrocyanation of non-conjugated ethylenically unsaturated compounds with nickel catalysts comprising phosphorous-containing ligands.\nAn integrated process for the production of ADN from BD and HC≡N can comprise BD hydrocyanation, 2M3BN isomerization to produce 3PN, and the hydrocyanation of pentenenitriles, including 3PN, to produce ADN and other dinitriles. Integrated processes are disclosed, for example, in United States Patent Application 2009/0099386 A1.\nDisclosed in United States Patent Publication No. 2007/0260086, is a process for the preparation of dinitriles with an aim to provide for the recovery of a catalyst formed by a mixture of mono- and bidentate ligands and to be able to reuse the catalyst thus recovered in the hydrocyanation and/or isomerization stages.\nUnited States Patent Publication No. 2008/0221351 discloses an integrated process for preparing ADN. A first process step includes hydrocyanating BD to produce 3PN over at least one zero-valent nickel catalyst. A second process step of the integrated process involves hydrocyanating 3PN to produce ADN over at least one zero-valent nickel catalyst and at least one Lewis acid. In this integrated process, at least one of the zero-valent nickel catalysts used in one of the process steps is transferred into the other process step.\nPhenolic compounds, such as phenol and cresols, may be present as a catalyst impurity in catalysts used to react BD with HCN or to isomerize 2M3BN. Phenolic compounds may be produced by hydrolysis of phosphorus-conating ligands. Phenolic compounds may react with catalysts used to react 3PN with HCN. Such reactions of phenolic compounds with catalysts may reduce the catalytic activity of the catalysts."}
{"text": "1. Field of the Invention\nThe present invention relates to compact antenna apparatuses such as those used for portable radio terminals.\n2. Discussion of the Background\nIn recent years, it has been important for portable radio terminals to be compact and thin for improved portability. As portable radio terminals have been made compact, however, the effectiveness of the antenna used therein may deteriorate. This tendency becomes especially prominent with a built-in antenna such as an inverted F antenna.\nFIG. 11 shows a conventional compact antenna apparatus in which a line-shaped inverted F antenna is disposed on a conductor plate. A line-shaped inverted F antenna 1101 is connected to a conductor plate 1102 by a short-circuited line 1103 and a power-feed line 1104. A portable radio terminals having an operational frequency of 800 MHz is used as an example. The radio terminals, is shown simply by a conductor plate. The longitudinal length L of the conductor plate 1102 is set to one fourth the wavelength of a signal having the specified operational frequency, in the example, one fourth of the wavelength of an 800 MHz signal.\nFIG. 12 is a graph showing the radiation efficiency of the line-shaped inverted F antenna 1101 shown in FIG. 11. The horizontal axis indicates the longitudinal length L of the conductor plate 1102 shown in FIG. 11. It is understood from this graph that the efficiency becomes -3 dB or less when the longitudinal length of the conductor plate becomes one fourth the wavelength or less. This indicates that about half the power supplied to the antenna 1101 is lost in the antenna 1101. This deterioration may occur in a compact antenna such as a line-shaped inverted F antenna as a result of the compact conductor plate.\nAs described above, as a portable radio terminals have been made compact, antenna performance may deteriorate. Especially with a compact antenna apparatus such as an inverted F antenna mounted on a conductor plate having a length of about one fourth the wavelength, the degree of deterioration is high and stable communication may be impeded."}
{"text": "Industrial control systems, like those used in the oil and gas production industry, frequently include one or more remote terminal units (RTUs) and/or flow computers as key components in an operating process unit of a control system (e.g., at a wellhead oil production site). RTUs are used to interface a control system host with field devices (e.g., valves, valve positioners, switches, sensors, transmitters, etc.) configured to perform control functions such as opening or closing valves and measuring process parameters. RTUs enable this interface by communicating commands from the host to the field devices and by communicating data sent by the field devices back to the host. Such communications may be implemented via any of analog, digital, or combined analog/digital buses using any desired communication media (e.g., hardwired, wireless, etc.) and protocols (e.g., Fieldbus, Profibus®, HART®, etc.). Additionally or alternatively, RTUs may act as standalone devices that implement process control and data archiving independent of commands provided by the host (and/or without connection to the host)."}
{"text": "Examples of communication terminals include smartphones, portable phones, and tablet terminals. In recent years, the communication terminal enables executing applications in addition to perform a communication.\nExamples of applications executed using the communication terminal include an application called a “native application”. The native application is executed by an Operating System (OS) of a communication terminal. Hence, the native application has a high dependence on the communication terminal.\nExamples of applications executed using a communication terminal also include an application called a “Web application”. The Web application is executed by a browser of a communication terminal. Hence, the Web application has a low dependence on the OS.\nExamples of the Web applications include a Hyper Text Markup Language (HTML) 5 application. The HTML5 application conforms to the HTML5 standard advocated by the World Wide Web Consortium (W3C).\nThe HTML5 uses an origin. According to the standard advocated by the W3C, the origin includes a protocol, a domain (host), and a port number. As a technology related to the origin, a technology has been proposed wherein access control for each part in an HTML document constituting a Web page is performed according to the origin of the part in the document (see, for example, patent document 1).\nPatent document 1: Japanese Laid-open Patent Publication No. 2008-299414"}
{"text": "Recently the use of enzymes, especially of microbial origin, has become more and more common. Enzymes are used in several industries including, for example, the starch industry, the dairy industry, and the detergent industry. It is well known in the detergent industry that the use of enzymes, particularly proteolytic enzymes, has created industrial hygiene concerns for detergent factory workers, particularly due to the health risks associated with dustiness of the available enzymes.\nSince the introduction of enzymes into the detergent business, many developments in the granulation and coating of enzymes have been offered by the industry. See for example the following patents relating to enzyme granulation:\nU.S. Pat. No. 4,106,991 describes an improved formation of enzyme granules by including within the composition undergoing granulation, finely divided cellulose fibers in an amount of 2-40% w/w based on the dry weight of the whole composition. In addition, this patent describes that waxy substances can be used to coat the particles of the granulate.\nU.S. Pat. No. 4,689,297 describes enzyme containing particles which comprise a particulate, water dispersible core which is 150-2,000 microns in its longest dimension, a uniform layer of enzyme around the core particle which amounts to 10%-35% by weight of the weight of the core particle, and a layer of macro-molecular, film-forming, water soluble or dispersible coating agent uniformly surrounding the enzyme layer wherein the combination of enzyme and coating agent is from 25-55% of the weight of the core particle. The core material described in this patent includes clay, a sugar crystal enclosed in layers of corn starch which is coated with a layer of dextrin, agglomerated potato starch, particulate salt, agglomerated trisodium citrate, pan crystallized NaCl flakes, bentonite granules or prills, granules containing bentonite, Kaolin and diatomaceous earth or sodium citrate crystals. The film forming material may be a fatty acid ester, an alkoxylated alcohol, a polyvinyl alcohol or an ethoxylated alkylphenol.\nU.S. Pat. No. 4,740,469 describes an enzyme granular composition consisting essentially of from 1-35% by weight of an enzyme and from 0.5-30% by weight of a synthetic fibrous material having an average length of from 100-500 micron and a fineness in the range of from 0.05-0.7 denier, with the balance being an extender or filler. The granular composition may further comprise a molten waxy material, such as polyethylene glycol, and optionally a colorant such as titanium dioxide.\nU.S. Pat. No. 5,254,283 describes a particulate material which has been coated with a continuous layer of a non-water soluble, warp size polymer. U.S. Pat. No. 5,324,649 describes enzyme-containing granules having a core, an enzyme layer and an outer coating layer. The enzyme layer and, optionally, the core and outer coating layer contain a vinyl polymer.\nWO 91/09941 describes an enzyme containing preparation whereby at least 50% of the enzymatic activity is present in the preparation as enzyme crystals. The preparation can be either a slurry or a granulate.\nWO 97/12958 discloses a microgranular enzyme composition. The granules are made by fluid-bed agglomeration which results in granules with numerous carrier or seed particles coated with enzyme and bound together by a binder.\nHowever, even in light of these developments offered by the industry (as described above) there is a continuing need for low-dust granules. In particular, it is especially problematic in the detergent industry when granules in general, or those comprising proteins or enzymes, form dust and are aerosolized. In these cases, workers are often exposed to the contents of the granules and can develop severe allergic reactions. Therefore, it is an object of the present invention to provide a method of producing a low-dust enzyme granule by adding antifoam agent. It is a further object of the invention to facilitate a safer environment for workers in the detergent industry who are exposed to enzyme containing granules."}
{"text": "Integrated Circuit (IC) packages, chips (sometimes called die), and other devices have evolved into very dense packaged structures. Semiconductor miniaturization has resulted in the development of very large scale integrated circuit (VLSI) devices with millions of active and passive components. These devices are typically encapsulated in a protective package providing a large number of pin-outs for mounting or interconnection to external circuitry through a carrier substrate such as a printed circuit board or other higher-level packaging.\nThe semiconductor industry is well over 60 billion dollars a year market. Improved fabrication technologies have lowered cost while increasing functional power thereby revolutionizing the electronic marketplace. The ability to put more and more functionality and higher performance is generally made possible by repeatedly shrinking feature sizes. However, as the more and more functionality and performance characteristics are included in these chips, the limits of current technologies and methodologies are approached. There are a number of issues and trends that are becoming readily apparent. For example, package temperature issues are exacerbated by a number of factors including static power and lowered junction temperatures due to smaller feature sizes and lowered thresholds. Also, hot spots within the chips are exacerbated by the active elements being buried so deeply beneath the conducting layers and contacts.\nIn the formation of integrated circuits, a number of active and passive semiconductor devices are formed on each of many die on a wafer, such as silicon. The fabrication technology for integrated circuits has vastly improved yields such that large arrays of electronic circuits on many die are produced on a single semiconductor wafer. In a typical die, two or more silicide layers and up to eight metallization layers, each separated by a dielectric layer, are currently used to interconnect among the active and passive elements to form the individual circuits and to interconnect among those circuits to form the die. They are also used to provide bonding pads at the top of the die. For example, Texas Instruments has a 65 nm process that uses eleven (11) copper conducting layers plus two polysilicide layers for a total of twenty-seven (27) conducting and dielectric layers.\nTypical interface schemes for integrated circuit packaging include Pin Grid Array (PGA), Ball Grid Array (BGA), and Land Grid Array (LGA). PGA packages use a two-dimensional array of pins directly connected by soldering or inserted into through-hole pads in a Printed Circuit Board (PCB). BGA packages have a two-dimensional array of conductive pads, such as balls, bumps or pillars, and are mounted by soldering the pads on the package to corresponding surface pads on the mount side of the PCB. LGA packages have an array of metal stubs and are mounted to the PCB in a clamp with a compressible interposer material placed between the package and the PCB.\nFor illustrative purposes, the description herein will focus somewhat on the most common package, the BGA. The BGA bond pads of the semiconductor die are sometimes connected to the printed circuit board via conductors, either by direct contact in a flip-chip orientation through conductive balls, bumps or pillars or, by intermediate connector elements comprising wire bonds, or TAB (flexible circuit) connections. More commonly they are connected inside a package which in turn is soldered to the PCB. In addition, as ball sizes shrink, they are less tolerant of the already worsening planarity problems.\nOne of the well-known problems with having the connection balls on the active surface of the die is the planarity required of the die. The problem is compounded as the number of mask levels increase to accommodate the latest demand for processors and memory. The further layers exacerbate the planarity problem and make the die connection more difficult. Another problem associated with the dense packaging is that the fanin/fanout of a given ball is limited to either a single fanin or a single fanout. One of many deleterious effects of this property is that only one element can be placed on a bus unless another delay stage is introduced. Typically, many three state elements are “hung” on a bus, and only one such element is active at any one time. The third state, the high impedance state, of the inactive elements ensures that they neither charge nor discharge the bus.\nWith respect to the formation of BGA structures, these structures typically require the die to be flipped upside down hence the name “flip chip”. A problem with flipping the chip upside down is that all possibilities of monitoring chip behavior in-situ are eliminated. Integrated circuit (IC) dies typically connect to the substrate within the IC package using either wire bond or Flip-Chip technology. Flip-Chip bonding is normally used for high pin count IC dies, and the pins on the Flip-Chip die are called bump pads. As with the package array technologies, there is a matching pattern of pads on the package substrate. Interconnect on the package substrate is typically used to connect the pads on the substrate (connected directly to the IC die) to the pins, pads, or stubs on the surface of the package that gets inserted, soldered, or pressed to the PCB.\nA further problem with using the current flip chip techniques arises from a combination of bringing the interconnects to the top of the die in combination with the large number of layers. The connection to/from the underlying layers must be brought to the top of the die for connection to a connection ball. Current chips have six to eight metallization layers plus two polysilicon layers, and each of these conducting layers is separated from adjacent conducting layers by a dielectric layer. Thus, a signal on the lowest conducting layer may have to go through as many as nine conducting layers and ten dielectric layers, including one on the top surface, in order to reach the top surface and be connected to an interconnect ball.\nIn order to traverse from each conduction layer through the intervening dielectric layer to connect to an adjacent conduction layer generally requires a stepped via to minimize and avoid step coverage problems. A stepped via is a hole whose cross-section is larger on one end than on the other. It has a contact on both conducting layers, and a short piece of interconnect to route to a via to go to the next upper or lower conduction layer. This process is repeated for each such signal line for each conducting layer until the top is reached. Both the vias and the contacts are considerably larger than the width of the line; thus a considerable area is utilized just for the interconnects to the balls on the top surface of the chip.\nIt is well-known that as the capacity and speed of many integrated circuit devices, such as dynamic random access memories (DRAMs) have increased, the number of inputs and outputs (I/Os) to/from each die has increased, requiring more numerous and complex external connections thereto and, in some instances, requiring undesirably long traces to place the bond pads serving as I/Os for the typical die in communication with the traces of the carrier substrate.\nA related problem with the current techniques is the interconnectivity between dies. Signals that travel to other dies must traverse not just upwards in one die, but may also travel further upwards if the second die is on top of the first die. The signals may also traverse downward into a second die to reach the desired point on a given conduction layer in the second die. Thus, for example, the signal may traverse upwards through many conduction and dielectric layers to reach the top die surface and go through a ball, a bonding pad, or other connection structure on the first die. After a connection to a second die, the signal would then be similarly routed downwards from the top surface of the second die to the desired layer in the second die. Thus, for example, a given signal could traverse about thirty-eight (fifty-four in the case of the 65 nm process of Texas Instruments) conduction and dielectric layers and their associated vias plus two conduction balls, bonding pads, or other connection structure. And, each of the contact layers adds resistance that, in combination with the total length traversed by the signals, may impact signal quality, speed, noise margin, bus skew, slack and setup/hold times, rise and fall times, potential spike generation and signal cancellation and makes timing convergence and other design and verification aspects difficult and sometimes limits upper clock frequencies.\nThus, such a tortured path impacts both speed and signal integrity. It is well known that interconnect delays are often much greater than gate delays below feature sizes of roughly 0.5 to 0.35 micron. Traversing substantial numbers of interconnect layers on two or more die considerably exacerbates the problems of timing convergence, layout, and architecture. Attempts to minimize these problems typically place restrictions on the layers regarding the manner in which signals can travel among these layers. These restrictions increase the already difficult layout and hinder designer creativity. They also sometimes force restrictions on the architecture of the chip.\nLong path lengths also add significantly to the capacitance seen by the source and slow the signal down. Moreover, as line widths continually decrease, 22 nanometers currently in development, the number of “squares” and hence the resistance of the line in a given length of line increases by about the same factor as the decrease in line widths. This then increases the resistance seen by the driving source, again by about the same factor as the decrease of the line widths, all else being the same.\nFurthermore, the increased line resistance adds to often already strained power usage and dissipation while making it more difficult to make bus signal elements and slave clock signals track their respective functions.\nIn addition, as higher speed IC assemblies operate at lower operational signal voltages, noise problems also become problematic. Mutual inductance results from an interaction between magnetic fields created by signal currents flowing to and from a packaged IC die through leads or traces, while self inductance results from the interaction of the foregoing fields with magnetic fields created by oppositely-directed currents flowing to and from ground. Signal propagation delays, switching noise, and crosstalk between signal conductors resulting from mutual and self inductance of the conductive paths contribute to signal degradation.\nWhile lead inductance in IC packages has not traditionally been troublesome, the increasing signal frequencies of state-of-the-art electronic systems have substantially increased the practical significance of package lead or trace inductance. For example, at such faster signal frequencies, performance of IC die using extended leads or traces for external electrical connection is slower than desirable because the inductance associated with the elongated conductive paths required slows changes in signal currents through the leads or traces, prolonging signal propagation. In addition, digital signals propagating along the leads or traces are spreading out causing the signal components and signals themselves to disperse. While mild dispersion merely widens the digital signals without detrimental effect, severe dispersion can make the digital signals unrecognizable. In addition, reflection signals propagating along the leads or traces as a result of impedance mismatches between the lead fingers and associated IC die or between the leads or traces and external circuitry, may distort normal signals propagating concurrently with the reflection signals. And, magnetic fields created by signal currents propagating, through the lead or trace-associated inductance can induce currents in adjacent leads or traces, causing crosstalk noise.\nTherefore, the state of the art die and package configurations described herein are having difficulties in keeping up with the trend towards faster devices at lower power. Noise problems are exacerbated by use of a large number of laterally adjacent traces of substantial and varying lengths extending from centralized die location to the horizontally-spaced, offset locations of vias extending to solder balls or other conductive elements for securing and electrically connecting the package to a carrier substrate.\nThere are various constraints that limit the number of signal traces that can be fabricated on a package. Industry standards impose specific requirements as to the spacing between solder balls, thereby restricting the spacing between the vias that electrically connect the signal traces to the solder bumps. The spacing restriction limits the number of signal traces that can fit between the vias which, in turn, limits the number of signal traces that can be used to carry signals to and from the die. Current fabrication technology imposes minimum pitch requirements for signal traces to attain satisfactory yields and to ensure mechanical and electrical reliability. The limitation on the maximum number of usable signal traces limits the maximum number of solder bumps, thereby placing a ceiling on the number of signals that a particular package can provide.\nThere are several levels of interconnections that make up an encapsulated integrated circuit. First, there are the internal interconnections of the circuit making up an individual die. Then there is the interconnection between the various dies themselves, especially if more than one die is in a package. There is also a connection layer between the dies and the package that allows external access.\nWith respect to the package connection, wire bonds are typically used to electrically connect the input/output (I/O) pads of the die to traces or pads on the package. If the die is on the side of the circuit board inside the package opposite the solder bumps, conductive vias are formed through the circuit to conduct signals from the solder bumps to the pads or traces. To enable routing in highly dense integrated circuit packages, there are many techniques known to those skilled in the art such as micro-vias, blind vias, buried vias, staggered vias.\nA further problem with higher density and more complex circuit design relates to the verification at all the various steps along the process. As the number of layers and the die size increase, the resources to validate the different steps including the design rules of the layouts of the different layers increase exponentially.\nIn addition, the static power is greater due to leakage through the thin oxides required for the higher density circuits. There is also an increase in the number of conduction and concomitant dielectric layers which must be traversed to reach various parts of a given die and to interact with other dies. This increase gives rise to increased problems with layer to layer registration, step coverage with its potential failure modes, contact integrity, and validation problems. Because of aspect ratios, stepped vias must generally be used in most cases at the smaller feature sizes. However, using the stepped vias adds an additional mask layer for each such stepped via. Planarity also becomes more difficult with the denser circuits.\nAs noted herein, the existing packaging fail to address several criteria required for high density, high speed and low power packaging, namely: interconnections with non-planar die resulting from many layers on the die; limitations on the fanin/fanout; inability to monitor flip-chip packaging signals; and, very long signal travel paths resulting in signal reflection, degraded signal integrity, propagation delays, switching noise, crosstalk and dispersion. Therefore, what is needed is a mechanically and electrically desirable packaging scheme to accommodate high density high speed, low power packaging. Such a system should interconnect die without having to travel excessive paths and allow in-situ monitoring of critical signals. Ideally the system should be easily implemented without affecting the current and expected fabrication methods and machinery."}
{"text": "Solid gas sorption systems are used to produce cooling and/or heating by repeatedly desorbing and absorbing the gas on a coordinative complex compound formed by absorbing a polar gas refrigerant on a metal salt in a sorption reaction sometimes referred to as chemisorption. Complex compounds incorporating ammonia as the polar gaseous refrigerant are especially advantageous because of their capacity for absorbing large amounts of the refrigerant, often up to 80% of the absorbent dry weight. The complex compounds also exhibit vapor pressure independent of the refrigerant concentration and can be made to absorb and desorb very rapidly. Apparatus using complex compounds to produce cooling are disclosed, for example, in U.S. Pat. Nos. 5,161,389, 5,186,020, and 5,271,239. Improvements in achieving high reaction rates for the complex compounds are achieved by restricting the volumetric expansion of the complex compound formed during the absorption reaction of the gas on the metal salt. The methods and apparatus for achieving such high reaction rates are disclosed in U.S. Pat. Nos. 5,298,231, 5,328,671, 5,384,101 and 5,441,716, the descriptions of which are incorporated herein by reference.\nWhile increased reaction rates have resulted from the aforesaid methods, it has been determined that repeated and relatively long-term absorption and desorption cycling of the complex compounds, particularly those using ammonia as a refrigerant, leads to sorbent migration even in the confined reaction chamber. It has also been found that the sorbent migration increases as higher sorption rates are used. Such migration may lead to uneven sorbent densities which in turn cause force imbalances in the heat exchanger structure, often resulting in deformation of the heat transfer surfaces and/or internal structures. As the heat exchanger structure becomes modified or compromised, heat and mass transfer reductions occur as does the sorption rate of the process. As sorbent migration continues, significant losses in performance efficiency are realized as is the possibility of failure of the reactor especially where it is exposed to high reaction rate sorptions.\nAlthough improvements in attempts to overcome sorbent migration have been made for metal hydrides, such procedures and techniques have not been found to be suitable for ammoniated complex compounds. In U.S. Pat. No. 4,507,263, there is described micro-immobilization for metal hydride using a sintering process in which a metal hydride powder is embedded in a finely divided metal and the mixture sintered in a furnace at 100-200.degree. C. using hydrogen pressure of 250-300 atmospheres. Although such a process reportedly results in mechanical stability for metal hydrides even after 6,000 cycles, the process is not effective for ammoniated complex compounds which exhibit much larger forces as compared to those experienced with metal hydrides. For example, where ammoniated complex compounds are absorbed and/or desorbed above about 3 moles NH.sub.3 /mole sorbent-hr, the forces exercised on a sintered metal structure are so large as to result in deformation of the structure. Moreover, for most practical applications using complex compound technology, practical life expectancy of the reactors will exceed 6,000 cycles by an order of magnitude."}
{"text": "1. Field\nExemplary embodiments of the present invention relate to a nonvolatile memory device and a method for fabricating the same, and more particularly, to a nonvolatile memory device with a three-dimensional structure in which a plurality of memory cells are vertically stacked from a substrate and a method for fabricating the same.\n2. Description of the Related Art\nA nonvolatile memory device is a memory device which maintains data stored therein even when power supply is interrupted. Currently, various nonvolatile memory devices, for example, a flash memory and the like are widely used.\nRecently, as improving the degree of integration of a nonvolatile memory device with a two-dimensional structure in which memory cells are formed in a single layer on a semiconductor substrate reaches a fabrication limit, a nonvolatile memory device with a three-dimensional structure in which a plurality of memory cells are formed along a channel layer vertically projecting from a semiconductor substrate has been suggested in the art. In detail, nonvolatile memory devices with a three-dimensional structure are generally divided into a structure which has a straight channel layer and a structure which has a U-shaped channel layer.\nIn the case of the conventional structure having the straight channel layer, while manufacturing processes are relatively simple and easy, since a source line is formed by implanting impurities into a silicon substrate, there are concerns in that a doping profile is likely to be changed by a subsequent annealing process and the like and source resistance is likely to increase. Although a method for forming a source line using a conductive substance such as a metal to reduce resistance has been proposed in the art, it may be difficult to control the doping profile of the source region under the channel layer.\nIn the case of the conventional structure having the U-shaped channel layer, while the above-described concerns may be solved, since the gate electrode of a pass transistor is formed long, channel resistance may increase. According to this fact, concerns are likely to occur in that the threshold voltage of the pass transistor increases and memory cell driving current is reduced."}
{"text": "The present invention relates to a light guide plate and a liquid crystal display (LCD) using the light guide plate.\nA liquid crystal display (LCD) is one kind of the most widely used flat panel displays due to its advantages such as compact volume, low power consumption, and no radiation. A LCD is formed by assembling an array substrate and a color filter substrate that are opposite to each other. A gate line and a data line intersect each other and are formed on the array substrate to define a pixel region in which a pixel electrode and a thin film transistor (TFT) as a switching element are disposed. A gate signal is applied over the gate line to a gate electrode of the TFT to turn on/off the TFT, and image data signal is applied to the pixel electrode over the data line. A black matrix is disposed on the color filter substrate to prevent light leakage in areas other than the pixel region. A color filter layer is disposed in each pixel region on the color filter area, and a common electrode is disposed on the color filter substrate. A liquid crystal (LC) layer is filled in a space between the array substrate and the color filter substrate to form a LCD panel. The LCD displays images by applying voltages to the pixel electrode and common electrode to generate an electric field in the LC layer, which determines orientations of LC molecules in the LC layer to change polarization of the incident light.\nA backlight module is an important component of a LCD and used for providing the LCD panel with light for display and for determining the color representation quality of the LCD. A conventional LCD backlight module is shown in FIG. 1. The backlight module comprises a light source 1 (for example, a cold cathode fluorescent lamp or light emitting diode), a reflective cover 2 around the light source 1, and a reflection plate 3, a light guide plate 4, a first diffuse plate 5, a first prism plate 6, a second prism plate 7, and a second diffuse plate 8 that are stacked in order. A portion of the light from the light source 1 directly enters the light guide plate 4, and another portion is reflected back by the reflective cover 2 and then enters the light guide plate 4. The light is reflected by the underlying the reflection plate 3 and exits from the upper surface of the light guide plate 4 to transmit through the overlaying optical sheets such as the first and second diffuse plates 5 and 8 and the first and second prism plates 6 and 7 before entering a LCD panel. The light guide plate 4 in the above backlight module serves as an optical medium. The overlaying optical sheets such as the first and second diffuse plates 5 and 8 and the first and second prism plates 6 and 7 are used to uniform the light and enlarge view angle."}
{"text": "This invention relates to a combination of strap and buckle suitable for diving fins.\nIn the accompanying drawings, FIG. 5 is a perspective view showing one of diving fins 101 disclosed in Japanese Patent Publication No. 1989-38510B as partially broken away and FIG. 6 is a sectional view taken along a line VIxe2x80x94VI in FIG. 5. It should be understood here that, in FIG. 6, the buckle 106 made of hard plastic material is shown in a sectional view and a heel strap 103 made of soft rubber is shown in a side view. This buckle 106 serves to connect the strap 103 to a rear end of a foot pocket 102 and to tighten the strap 103 around a fin wearer\"\"s heel. The buckle 106 comprises a sheath-shaped first coupling member 104 and a second coupling member 105 wherein the first coupling member 104 is mounted on the rear end of the foot pocket 102. The second coupling member 105 has a insert portion 114 releasably attached to the first coupling member 104 and a holding frame mechanism 115 serving to hold the strap 103 draped around a pin 119 provided in the holding frame mechanism 115 as its component. This holding frame mechanism 115 includes movable locking means 120 supported by a shaft 123 and biased by a spring 128 so as to be rotated clockwise as viewed in FIG. 6. The second coupling member 105 is formed on its lower surface with a second locking tooth 124. The strap 103 draped around the pin 119 of the holding frame mechanism 115 is formed on its one surface with a plurality of first locking teeth 121 arranged in longitudinal direction of the strap 103 so that one of these first locking teeth 121 may bear against the second locking tooth 124 in a direction indicated by an arrow 129 in FIG. 6. The second locking tooth 124 is moved off upward from the first locking tooth 121 as the movable locking means 120 is rotated against the biasing effect of the spring 128. Thereupon, the strap 103 may be pulled in the direction indicated by the arrow 129 to move the strap 103 in a direction in which the strap 103 is slackened. Each of the first locking teeth 121 has an inclined rear surface 121A so that the rear surface 121A causes the movable locking means 120 to float as the strap 103 is pulled in the direction opposite to the direction indicated by the arrow 129. In this way, the strap 103 can be moved in the direction in which the strap 103 is tightened.\nIn the convention combination of buckle 106 and strap 103 as shown in FIGS. 5 and 6, the rectilinear strap 103 is draped around the pin 119 so as to be forcibly flexed in U- or V-shape. The strap 103 normally tends to restore its initial rectilinear state, i.e., tends to move in a direction indicated by an arrow 140 in FIG. 6. Consequently, one of the first locking teeth 121 bears against a rear end 115A of the holding frame mechanism 115 from the right as viewed in FIG. 6 as the strap 103 is pulled in the direction indicated by the arrow 129 and often prevents the strap 103 from smoothly moving in the direction indicated by the arrow 129. Thus, operation of slackening the strap 103 would take much time.\nIt is an object of this invention to improve the conventional combination of buckle and strap so that the strap can be smoothly moved and slackened without taking much time and labor.\nAccording to this invention, there is provided a combination of buckle and strap for diving fins adapted to couple a heel strap to a foot pocket and at the same time to adjustably tighten the strap around a diver\"\"s heel.\nThe strap is formed on its one surface a plurality of first locking teeth arranged intermittently in a longitudinal direction of the strap, each of the first locking teeth extending in a transverse direction of the strap so as to leave transversely opposite side edge portions of the strap free from any tooth. The buckle comprises a sheath-shaped first coupling member and a second coupling member releasably engaged with the first coupling member and having front and rear ends, wherein the second coupling member comprises an insert portion formed adjacent the front end so as to be releasably engaged with the first coupling member and a holding frame mechanism formed adjacent the rear end within which the strap is draped around so as to be movable in the longitudinal direction of the strap. The holding frame mechanism includes a movable locking piece pivotally supported by the holding frame mechanism. The movable locking piece has a second locking tooth normally biased by springs to bear against the first locking tooth and thereby to prevent the strap from moving in a direction in which the strap is slackened and pivotally moved upward away from the first locking tooth against the biasing effect of the springs, allowing the strap to move in a direction in which the strap is tightened, as the strap is pulled in the direction in which the strap is tightened. In the vicinity of the rear end, the holding frame mechanism presents an inverted U-shaped section defined by a top portion and a pair of lateral portions depending from transversely opposite side edges of the top portion wherein a pin around which the strap is draped extends between the pair of lateral portions below the top portion and a lower surface of the top portion has transversely opposite side edge portions bears against the strap disposed between the lower surface and the pin so that the transversely opposite side edge portions of the strap extending outside the first locking teeth may slidably move and transversely middle zone of the lower surface of the top portion is hollowed out upward by a dimension deeper than a height of the first locking teeth so that the first locking teeth can move in the longitudinal direction of the strap."}
{"text": "The present invention relates generally to plastic fasteners and more particularly to clips of plastic fasteners.\nPlastic fasteners are well known in the art and are commonly used to couple together two or more objects. For example, plastic fasteners have been used to attach merchandise tags to articles of commerce, to couple or to re-couple a button to an article of clothing, to last together shoe uppers, and the like.\nPlastic fasteners typically comprise a flexible filament having a first enlargement at one end thereof and a second enlargement at the opposite end thereof. In one common type of plastic fastener (see, for example, FIG. 1 of U.S. Pat. No. 5,321,872, which patent is incorporated herein by reference), the first enlargement has the shape of a transverse bar and the second enlargement has the shape of a paddle or the shape of a second transverse bar, the transverse bar and the paddle (or second transverse bar) extending in planes parallel to one another. In another common type of plastic fastener (see, for example, U.S. Pat. No. 3,494,004, which patent is incorporated by reference), the first enlargement has the shape of a transverse bar and the second enlargement has the shape of a knob or pin head. In still another common type of plastic fastener (see, for example, U.S. Pat. No. 4,240,183, which patent is incorporated herein by reference), the first enlargement has the shape of a transverse bar or the shape of a plug and the second enlargement has the shape of a socket, said socket being adapted to receive said transverse bar or said plug.\nPlastic fasteners of the various types described above are typically molded as parts of a unitary fastener clip. An example of such a fastener clip is disclosed in U.S. Pat. No. 3,733,657, which patent is incorporated herein by reference. The clip of the aforementioned '657 patent includes a plurality of fasteners, each of said fasteners comprising a flexible filament having a transverse bar (or \"cross-bar\") at one end thereof and a paddle or a second transverse bar (or \"cross-bar\") at the opposite end thereof, the transverse bar and the paddle (or second transverse bar) of each fastener extending in planes parallel to one another. The fasteners are arranged relative to one another so that the respective transverse bars are spaced apart and oriented side-by-side and parallel to one another and so that the respective paddles (or second transverse bars) are spaced apart and oriented side-by-side and parallel to one another. The clip of the foregoing '657 patent also includes a runner bar, said runner bar extending perpendicularly relative to the respective transverse bars and being connected to each of the transverse bars by a severable connector. The clip of said '657 patent further includes a severable member interconnecting each pair of adjacent paddles (or second transverse bars).\nOther examples of fastener clips that comprise one or more runner bars interconnecting plastic fasteners include U.S. Pat. Nos. 5,622,257, 4,901,854, 5,799,375, and U.S. Reissue Patent No. 34,891, all of which are incorporated herein by reference.\nFastener clips which comprise a runner bar suffer from certain disadvantages. For example, because the runner bar of a fastener clip is of no use once the fasteners originally attached thereto have been dispensed therefrom, a used runner bar typically constitutes economically and environmentally undesirable waste material. In addition, severed connectors previously used to connect fasteners to a runner bar and still remaining on the runner bar after the fasteners have been detached therefrom often have an acute end which may undesirably snag on and damage merchandise when fasteners from the fastener clip are dispensed into such merchandise with a conventional fastener dispensing tool. Furthermore, the number of fasteners that can be molded into a clip of the type having a runner bar is typically limited by the molding process used to create the fastener clip.\nFor at least the above reasons, a number of runner bar-less fastener clips, assemblies or stock have been fashioned. For example, in U.S. Pat. No. 4,039,078, inventor Bone, which issued Aug. 2, 1977, and which is incorporated herein by reference, there is disclosed fastener attachment stock to be separated or divided, e.g., by cutting, severing, rupturing or shearing, to provide a plurality of fasteners each preferably having a substantially H shape. The stock in its most preferred form includes two undivided elongated and continuous plastic side members having a plurality of plastic cross links coupled to and between each of said side members, each of the links being preferably spaced equidistantly apart from each other.\nFastener stock related to that disclosed in U.S. Pat. No. 4,039,078 is disclosed in U.S. Pat. No. 4,456,123, inventor Russell, which issued Jun. 26, 1984, and which is incorporated herein by reference. The fastener stock of U.S. Pat. No. 4,456,123 differs from that of U.S. Pat. No. 4,039,078 in that, in U.S. Pat. No. 4,456,123, the filament has a substantially D-shaped cross-section and provides draft on surfaces extending from the plane to facilitate removal from the mold. Also, the side members are reduced in crosssectional area between individual fasteners to provide severable connectors to facilitate separation. The connectors join the end-bars of adjacent fasteners end-to-end at a portion of their periphery, preferably having a flat face at said plane extending from said plane on either the same side as the filaments or the opposite side thereof. Preferably, they extend from the same side and the joined end-bars are substantially D-shaped. Where the connectors extend from the opposite side, the section thereof is preferably continued across the joined end-bars to provide a more rounded cross-section for feeding through circular needle bores.\nAs can be seen, in each of U.S. Pat. Nos. 4,039,078 and 4,456,123, the enlargements of adjacent fasteners are oriented in an end-to-end relationship. In contrast, in U.S. Pat. No. 4,660,718, inventors Kato et al., which issued Apr. 28, 1987, and which is incorporated herein by reference, there is disclosed a runner bar-less fastener assembly comprising fasteners of the type comprising a flexible filament having a head at one end thereof and a cross-bar at the opposite end thereof, wherein the respective heads and cross-bars of adjacent fasteners are oriented in a parallel, side-by-side, spaced relationship. More specifically, the two side faces of each head are formed so as to protrude or bulge from edges towards a central portion to provide an apex, and adjacent heads are mutually and directly connected through their facing apices in a manner capable of being easily disconnected by cutting. The sides of adjacent cross-bars are connected by a film or a rod that extends longitudinally along a substantial portion of the length of the cross-bars or by a plurality of connectors posts that extend transversely relative to the sides of adjacent cross-bars. As explained in the foregoing '718 patent, the purpose of the aforementioned film, rod and connector posts is to prevent adjacent cross-bars from moving, i.e., pivoting, relative to one another.\nIn U.S. Reissue Patent 32,332, inventor Kato, which reissued Jan. 20, 1987, there is disclosed a runner bar-less fastener assembly comprising fasteners of the type comprising a flexible filament having a head at one end thereof and a cross-bar at the opposite end thereof, wherein the respective heads and cross-bars of adjacent fasteners are oriented in a parallel, side-by-side, spaced relationship. More specifically, the two side faces of each head are formed so as to protrude or bulge from edges towards a central portion to provide an apex, and adjacent heads are mutually and directly connected through their facing apices in a manner capable of being easily disconnected by cutting. Likewise, the two side faces of each cross-bar are formed so as to protrude at a central portion, and adjacent cross-bars are directly connected through their respective protruding portions in another embodiment, the aforementioned directly-conjoined heads are replaced with a second set of cross-bars that are directly conjoined in the same manner as described above. According to the foregoing '332 patent, the direct conjoining of adjacent heads to one another and the direct conjoining of adjacent cross-bars to one another, both in the manner described above, is preferable to the use of connector posts for the reason that, when severed, connector posts are said to leave unwanted whisker-like projections on the heads and/or cross-bars to which they are attached."}
{"text": "Supply chain planning, which comprises the logistical plan of an in-house supply chain, is essential to the success of many of today's manufacturing firms. Most manufacturing firms rely on supply chain planning in some form to ensure the timely delivery of products in response to customer demands. Typically, supply chain planning is hierarchical in nature, extending from distribution and production planning driven by customer orders, to materials and capacity requirements planning, to shop floor scheduling, manufacturing execution, and deployment of products. Supply chain planning ensures the smooth functioning of different aspects of production, from the ready supply of components to meet production demands to the timely transportation of finished goods from the factory to the customer.\nA modern supply chain often encompasses a vast array of data. The planning applications that create and dynamically revise plans in the supply chain in response to changing demands and capacity require rapid access to data concerning the flow of materials through the supply chain. The efficient operation of the supply chain depends upon the ability of the various plans to adjust to changes, and the way in which the required data is stored determines the ease with which it can be accessed.\nOperations in a production line of a plant for producing a product are carried out at work centers. In supply chain planning tools work centers are represented by business objects that can e.g. represent the following real work centers: machines or machine groups; production lines; assembly work centers; and employees or groups of employees.\nTogether with bills of material and routings, business objects representing work centers belong to the master data in production planning and control systems. Business objects representing work centers are used in task list operations and work orders. Task lists are for example routings, maintenance task lists, inspection plans and standard networks. Work orders are created for production, quality assurance, plant maintenance and for the project system as networks.\nData in work centers is used for scheduling, costing, capacity planning, and simplifying operation maintenance. For the purpose of scheduling, operating times and formulas are entered in the business object representing the work center, so that the duration of an operation can be calculated. For the purpose of costing, formulas are entered in the business object representing the work center, so that the costs of the operation can be calculated. Usually a business object representing the work center is also assigned to a cost center.\nThe available capacity and formulas for calculating capacity requirements are entered into the business object representing a work center for capacity planning. Further, various default values for operations can be entered in the business object representing the work center for simplifying operation maintenance.\nA business object representing a work center is created for a plant and is identified by a key. The work center category that can be defined in customising, determines which data can be maintained in the business object representing the work center.\nSupply chain planning and management tools as the SAP R/3 system use routing. Routing is a description of which operations, e.g. process steps, have to be carried out and in which order to produce a product. In addition to information about the operations and the order in which they are carried out, routing also contains details about the work centers at which they are carried out as well as about the required production resources and tools. Standard values for the execution of individual operations are also saved in routings. Usually a bill of material (BOM) is assigned to routing. Individual components of the BOM are assigned to the routing operations."}
{"text": "The present disclosure relates to an image forming apparatus.\nRecently, image forming apparatuses have a function to be set to a power saving state in which for example power sources for some devices included therein are turned off (also referred to below as a sleep operation mode) in order to reduce power consumption. However, in a situation in which a user uses such an image forming apparatus in the sleep operation mode, it takes a waiting time to return the mode of the image forming apparatus from the sleep operation mode to a normal operation mode through touch panel input or key input. A technique using a motion sensor as described below is proposed to reduce the waiting time.\nAn image forming apparatus stores a detection state of the motion sensor and a state of a user using the image forming apparatus for a specific period and changes sensitivity of the motion sensor. When the motion sensor is used, mis-detection resulting in reversion to the normal operation mode from the sleep operation mode may occur due to a person who does not use the image forming apparatus approaching the image forming apparatus, besides detection resulting in normal reversion to the normal operation mode from the sleep operation mode due to a user approaching the image forming apparatus. A technique to change the sensitivity of the motion sensor is proposed to solve the above problem."}
{"text": "1. Field of the Invention\nThe present invention relates to electrical devices, and more particularly to methods and circuits for programming of antifuses.\n2. Description of Related Art\nProgrammable semiconductor devices include programmable read only memories (\"PROMs\"), programmable logic devices (\"PLDs\"), and programmable gate arrays. Programmable elements suitable for one or more of these device types include fuses and antifuses.\nA fuse is a structure which electrically couples its first electrode to its second electrode, but which, when programmed by passage of sufficient current between its electrodes, electrically decouples the first electrode from the second electrode.\nAn antifuse is a structure which when unprogrammed does not electrically couple its first and second electrodes, but which, when programmed, permanently electrically couples the first and second electrodes. An antifuse is programmed by applying sufficient voltage (\"programming voltage\") between its first and second electrodes. One type of antifuse comprises a high resistivity material in which a low resistivity filament is formed when the material is heated by electrical current. Amorphous silicon, silicon dioxide and silicon nitride have been used successfully as the high resistivity materials. See, for example, U.S. Pat. No. 4,823,181 issued Apr. 18, 1989 to Mohsen et al.; B. Cook et al., \"Amorphous Silicon Antifuse Technology for Bipolar PROMs,\" 1986 Bipolar Circuits and Technology Meeting, pages 99-100.\nAn antifuse, when programmed, should have a low resistance. It was generally believed that in order to obtain lower resistance one needs to raise \"programming\" current (the current passing through the antifuse during programming). Namely, the physics of antifuse programming was believed to be as follows. When the programming voltage is applied between the antifuse terminals, the high resistivity material breaks down at its weakest portion. Current flows through that portion and heats the material. The heat creates a conductive filament through the material. As the filament grows in size, the resistance across the material decreases. Hence the temperature of the material also decreases. Gradually the temperature becomes so low that the conductive filament stops growing. See Hamdy et al., \"Dielectric Based Antifuse for Logic and Memory ICs,\" IEDM 1988, pages 786-789. In order to reduce the resistance further, the current has to be increased so as to generate more heat.\nIt was confirmed experimentally that a higher programming current does provide a lower resistance. However, the current in a programmable circuit cannot be increased indefinitely because high current can damage circuit devices. Thus, it is desirable to find a programming method that provides a low antifuse resistance while using a low programming current.\nFurther, the resistance of the programmed antifuse varies from one antifuse to another even among antifuses of the same general construction, and even when the same technique is used to program the antifuses. Since the resistance is variable, designers and users of circuits with antifuses have to accommodate a wide range of antifuse resistances. There is a need for a programming method that would make the resistance less variable.\nA typical programmable circuit (for example, a gate array) contains hundreds or thousands of antifuses. The programming circuit must address the antifuses being programmed so as to program only those antifuses. At present, addressing circuits typically require a decoder. The decoder makes the programming circuit more complex. See, for example, U.S. Pat. No. 4,873,459 issued Oct. 10, 1989 to El Gamal et al.\nIt is desirable to provide a simpler programming circuit suitable for use in programmable circuits with many antifuses. In addition, such a programming circuit should consume little power. Further, the programming circuit should program a large number of antifuses fast."}
{"text": "Work machine operators can experience significant levels of vibration. Many regulatory bodies have imposed restrictions on the vibration levels that an operator may be exposed to over time. To comply with these restrictions, an operator's time on a particular machine can be limited. Specifically, the operator may be required to cease operation of the machine once he has experienced a certain vibration level for a predetermined period of time. Alternatively, an active vibration management system may be employed in an attempt to reduce the average vibration level experienced by the operator and, therefore, prolong his allowed time on the machine.\nVarious systems have been proposed for actively reducing vibrations in a machine. Many of these systems involve sensing of vibrations produced in the machine and reducing the vibrations transferred from a vibration source to the frame of the machine. For example, U.S. Pat. No. 6,644,590 to Terpay et al. (“the '590 patent”), which issued on Nov. 11, 2003, describes an active system and method for reducing vibrations generated by a gearbox in a rotary wing aircraft. In this system, an active mount is connected between the gearbox and the airframe using hydraulic actuators to suspend the airframe from the gearbox. Based on output signals from various vibration sensors, hydraulic fluid may be supplied to the actuators to move the gearbox relative to the airframe. This motion may be controlled to minimize the transfer of vibrations from the gearbox to the frame.\nWhile the system of the '590 patent may help reduce the vibrations transferred to certain machine components, the system has several shortcomings. For example, the system of the '590 patent cannot monitor or track average vibration levels experienced by an operator or component. Further, the system includes no predictive capability for determining the vibration response of a system to various operator inputs. In addition, the system does not include the capability of adjusting the response of a machine component to reduce the amount of vibration produced. Therefore, the system of the '590 patent may be unsuitable as a means for ensuring that an operator of a work machine does not experience a certain vibration level for greater than a permissible length of time.\nThe present disclosure is directed to overcoming one or more of the problems associated with the prior art active vibration reduction systems."}
{"text": "1. Field of the Invention\nThe present invention relates to a structure of a complementary metal insulator semiconductor device (CMIS) using a metal gate electrode, and a manufacturing method thereof.\n2. Description of the Related Art\nIn order to meet the requirements for an effective reduction in the thickness of a gate insulating film associated with increasing performance and increasing integration of semiconductor devices, it will be necessary in the future to introduce the technique for a metal gate electrode and a high dielectric (high-k) gate insulating film. To obtain proper performance in a CMIS transistor (cMISFET) using the metal/high-k gate insulating film, an effective work function φeff of a metal gate material needs to be about 3.9 to 4.3 eV for an n-channel MIS transistor (nMISFET), and about 4.8 to 5.2 eV for a p-channel MIS transistor (pMISFET).\nHowever, a metal having a low work function suitable for the n-channel MIS transistor is generally not stable in a heat treatment step necessary for a transistor formation process, and cannot have a φeff of about 3.9 to 4.3 eV suitable for the n-channel MIS transistor especially on the high-k gate insulating film after the formation of the transistor. Therefore, insertion of a layer containing groups IIA and IIIA metallic elements into a gate stack structure is necessary, which is effective as a technique of reducing Vth in the n-channel MISFET.\nOn the other hand, the layer containing the groups IIA and IIIA metallic elements increases Vth in the p-channel MIS transistor, so that there is a need for a step of detaching the layer containing the groups IIA and IIIA metallic elements in the p-channel MISFET region.\nHowever, the layer containing the groups IIA and IIIA metallic elements is generally low in resistance to an etching solution (e.g., refer to H. Y. Yu et al., Tech. VLSI, P18(2007)). Accordingly, there has been concern that the layer containing the groups IIA and IIIA metallic elements in the n-channel MISFET region may also be detached in the step of detaching the layer containing the groups IIA and IIIA metallic elements in the p-channel MISFET region or in an associated mask detaching step, and proper Vth modulation may not be obtained in the n-channel MISFET region.\nThus, there has been desired the realization of a CMIS structure which inhibits the effects of the Vth modulation by the layer containing the groups IIA and IIIA metallic elements in the p-channel MISFET region without performing the step of detaching the layer containing the groups IIA and IIIA metallic elements formed in the p-channel MISFET region."}
{"text": "The semiconductor integrated circuit (IC) industry has experienced rapid growth. Continuing advances in semiconductor manufacturing processes have resulted in semiconductor devices with finer features and/or higher degrees of integration. Functional density (i.e., the number of interconnected devices per chip area) has generally increased while feature size (i.e., the smallest component that can be created using a fabrication process) has decreased. This scaling-down process generally provides benefits by increasing production efficiency and lowering associated costs.\nDespite groundbreaking advances in materials and fabrication, scaling planar devices such as a metal-oxide-semiconductor field effect transistor (MOSFET) device has proven challenging. To overcome these challenges, circuit designers look to novel structures to deliver improved performance, which has resulted in the development of three-dimensional designs, such as fin-like field effect transistors (FinFETs). The FinFET is fabricated with a thin vertical “fin” (or fin structure) extending up from a substrate. The channel of the FinFET is formed in this vertical fin. A gate is provided over the fin to allow the gate to control the channel from multiple sides. Advantages of the FinFET may include a reduction of the short channel effect, reduced leakage, and higher current flow.\nHowever, since feature sizes continue to decrease, fabrication processes continue to become more difficult to perform. Therefore, it is a challenge to form a reliable semiconductor device including the FinFET."}
{"text": "With recent process miniaturization, the number of processing steps necessary to manufacture a nonvolatile memory such as a flash memory has been increasing. Increases in the numbers of times of deposition, etching, and exposure cause more defects on a semiconductor substrate, and therefore it becomes difficult to ensure reliability of a capacitance element including a gate oxide film on semiconductor substrate.\nFor example, JP-A-Heisei 2-76251 (Patent Literature 1) discloses a technique that automatically prevents defective operation due to insulation breakdown in a device including a large scale capacitance unit such as a charge pump circuit. A polysilicon electrode PSi divided into a plurality of sections is present on the surface of a semiconductor substrate, and an oxide film is present between the polysilicon electrode sections PSi and the semiconductor substrate, to form a capacitance. The polysilicon electrode sections are connected with narrow aluminum wiring lines Al.\nIn an ordinary state, the polysilicon electrode sections are all connected with a same wiring line, and therefore the plurality of electrode sections can operate as an integrated capacitance element. If an oxide film for a capacitance is broken through a specific polysilicon electrode portion in a state that a voltage difference is supplied between the substrate and the electrode, leakage current flows, concentrating on a broken electrode portion. In this case, a wiring line for connection between the electrode sections is narrow, and therefore the wiring line connected to the broken electrode section is fused and disconnected by an extraordinary leakage current due to the insulation breakdown to separate the broken polysilicon electrode section from the device.\nIn this technique, a fuse function for removing the defective device is covered by the wiring line. However, in the recent miniaturized process, to stably fabricate a narrow wiring line as a process node serving as the fuse is difficult in terms of a fabrication method. At a miniaturized process node, influence of a variation on a wiring line width is large, and therefore it is difficult to equalize a current value at which the fuse disconnects the connection to an arbitrary value.\nAlso, JP2003-338553A (Patent Literature 2) discloses a technique related to a semiconductor device provided with a nonvolatile memory that interrupts a leakage current to a broken electrostatic protective component."}
{"text": "Immunological assays are receiving attention as a method of identifying and quantitating, with high accuracy, proteins such as virus, bacteria, and allergenic substance contained in biological samples (for example, blood).\nThe immunological assay is roughly classified into an immunonephelometry and an immunolabeling assay. The immunonephelometry is a method of determining the change in turbidity of a sample solution that is caused by an antigen-antibody complex produced by an antigen-antibody reaction. The immunolabeling assay is a method of determining the change in the amount of a labeling substance after an antigen-antibody reaction by using an antibody labeled with the labeling substance.\nThe immunolabeling assay is subdivided according to the type of the labeling substance. Examples thereof include a radioimmunoassay in which a radioisotope is used as the labeling substance, an enzyme immunoassay (EIA) in which an enzyme is used, and a fluorescence immunoassay in which a fluorescent substance is used. In the EIA, as compared to other immunolabeling assays, the safety of the labeling substance is higher, it can be carried out by a simpler operation, and the measurement accuracy is higher. Therefore the EIA is used frequently.\nA typical example of the EIA is an enzyme-linked immunosorbent assay (ELISA). An example of the ELISA is described with reference to FIG. 1.\n<Step 1-A: Forming a Solid-Phase Antigen>\nA solution containing a capture antigen having an epitope identical to that of an antigen to be measured is introduced into a reaction chamber 7 and is maintained at a predetermined temperature for a predetermined period of time, and thereby the capture antigen is allowed to adsorb to the surface of the reaction chamber 7. Thereafter the capture antigen to be allowed to adsorb is covered with protein that is not involved in a later antigen-antibody reaction and enzyme reaction (blocking). Thus a solid-phase antigen 6 is formed inside the reaction chamber 7. The solid-phase antigen 6 is being fixed to the surface of the reaction chamber 7 and therefore is not removed from the reaction chamber 7 by the washing described later.\n<Step 1-B: Binding Reactions Between Antibody and Antigen to be Measured as Well as Solid-Phase Antigen>\nAn antibody 3 that specifically binds to an antigen to be measured 2 is added to a sample solution 9. The antibody 3 is being labeled with a labeling substance (an enzyme) 4. Thereafter, the sample solution 9 containing the antibody 3 is introduced into the reaction chamber 7. Thus, the antigen-antibody reaction proceeds between the antibody 3 and the antigen 2 to be measured as well as the solid-phase antigen 6 in the reaction chamber 7.\n<Step 1-C: Removal of Unreacted Antibody and Antigen-Antibody Complex>\nUsing a wash solution, the inside of the reaction chamber 7 is washed. Thereby the antigen-antibody complex formed through binding of the antibody 3 to the antigen 2 to be measured and the antibody 3 that has not been bound to the solid-phase antigen 6 are removed from the reaction chamber 7. Accordingly, a conjugate of the solid-phase antigen 6 and the antibody 3 remains in the reaction chamber 7.\n<Step 1-D: Measurement of the Amount of the Labeling Substance>\nThe amount of the labeling substance 4 of the conjugate of the solid-phase antigen 6 and the antibody 3 that has remained in the reaction chamber 7 is measured. This measurement is carried out, for example, as follows. First, a solution containing a measuring reagent (for example, a substrate of the enzyme) that reacts with the labeling substance 4 is prepared. The solution is prepared by adding the measuring reagent to a buffer solution typified by Tris-HCl buffer. Next, this solution is introduced into the reaction chamber 7, and thereby a measurement solution 10 is obtained that contains the measuring reagent and the conjugate of the solid-phase antigen 6 and the antibody 3. After the reaction between the measuring reagent and the labeling substance 4 is allowed to proceed in the measurement solution 10, a signal that reflects the amount of reaction product is detected.\nAs a result of this measurement, the amount of the antigen 2 to be measured in the sample solution 9 is calculated based on the amount of the solid-phase antigen 6 and the amount of the sample solution 9 introduced into the reaction chamber 7 in Step 1-B.\nFrom the viewpoint of measuring the amount of the substance to be measured in the sample solution using a trace amount of sample solution in a short period of time, a chip-type biosensor is receiving attention. For example, JP 2 (1990)-062952 A discloses a chip-type biosensor including an insulating chip substrate, an electrode system disposed on the chip substrate, an enzyme reaction layer disposed on the electrode system, and an insulating layer that has notches and is disposed above the chip substrate in such a manner that the electrode system and the enzyme reaction layer are exposed. The enzyme reaction layer contains a measuring reagent for inducing an enzymatic cycling reaction, which is typified by an oxidoreductase and an electron mediator. An enzyme containing as a substrate the substance to be measured is used as the oxidoreductase. For example, glucose oxidase is used when the glucose amount is to be measured, and cholesterol oxidase is used when the cholesterol amount is to be measured. In this biosensor, a sample solution is dripped into the notches, so that the measuring reagent is dissolved in the sample solution. Thereby the reaction between the enzyme and the substance to be measured through an electron mediator (an enzymatic cycling reaction) proceeds. The amount of the substance to be measured in the sample solution is calculated based on the oxidation current value that is obtained by electrochemically oxidizing the electron mediator reduced by the enzyme reaction.\nIn the conventional immunological assay as shown in FIG. 1, as described above, Step 1-D is carried out using the measurement solution 10 containing a buffer solution as a solvent. Accordingly, in order to carry out the immunological assay on a chip, it is necessary to supply the buffer solution from the outside of the chip or to allow the buffer solution to be retained on the chip beforehand. However, when the buffer solution is supplied from the outside of the chip, extra time and effort is required for the supply. When the buffer solution is allowed to be retained on the chip, this solution desirably is allowed to be retained on the chip in a hermetic state and in the state that facilitates it to be introduced into the reaction chamber for assaying while containing a component that reacts with the labeling substance. However, it is not easy to form such a retaining state on a chip."}
{"text": "Electronic mail (email) has become an integral part of business and personal communications. As such, many users have multiple email accounts for work and home use. Moreover, with the increased availability of mobile cellular and wireless local area network (LAN) devices that can send and receive emails, many users wirelessly access emails from mailboxes stored on different email storage servers (e.g., corporate email storage server, Yahoo, Hotmail, AOL, etc.).\nYet, email distribution and synchronization across multiple mailboxes and over wireless networks can be quite challenging, particularly when this is done on a large scale for numerous users. For example, different email accounts may be configured differently and with non-uniform access criteria. Moreover, as emails are received at the wireless communications device, copies of the emails may still be present in the original mailboxes, which can make it difficult for users to keep their email organized.\nOne particularly advantageous “push” type email distribution and synchronization system is disclosed in U.S. Pat. No. 6,779,019 to Mousseau et al., which is assigned to the present Assignee and is hereby incorporated herein by reference. This system pushes user-selected data items from a host system to a user's mobile wireless communications device upon detecting the occurrence of one or more user-defined event triggers. The user may then move (or file) the data items to a particular folder within a folder hierarchy stored in the mobile wireless communications device, or may execute some other system operation on a data item. Software operating at the device and the host system then synchronizes the folder hierarchy of the device with a folder hierarchy of the host system, and any actions executed on the data items at the device are then automatically replicated on the same data items stored at the host system, thus eliminating the need for the user to manually replicate actions at the host system that have been executed at the mobile wireless communications device.\nThe foregoing system advantageously provides great convenience to users of wireless email communication devices for organizing and managing their email messages. Yet, further convenience and efficiency features may be desired in email distribution and synchronization systems as email usage continues to grow in popularity. For example, in new user accounts, an email provisioning and authentication system can run through a series of possible email server configurations to determine how to access an electronic mailbox for a user email account. The user can supply email address parameters such as an email address and password, but often becomes frustrated if the wrong email address parameter is typed. The user typically must wait a relatively long time to determine if something is wrong, or worse, the user may be given an advanced configuration screen and asked to provide difficult to know IP address numbers, ports and other entries because of the mistake. Some prior art systems have parsed emails and tried to provision, and as a subsequent step after failure, used MX records to aid in the process for accessing email. But those systems have not been used for provisioning in a more direct manner."}
{"text": "This invention relates in general to a magnetically enhanced wafer processing system and relates more particularly to a wafer processing plasma reactor in which the processing rate at the wafer is adjusted by means of magnetic material included within the reactor processing chamber.\nPlasma processing, such as plasma etching and deposition processing, in the fabrication of circuits is attractive because it can be anisotropic, can be chemically selective and can produce processing under conditions far from thermodynamic equilibrium. Anisotropic processing enables the production of integrated circuit features having sidewalls that extend substantially vertically from the edges of a the masking layer. This is important in present and future ULSI devices in which the depth of etch and feature size and spacing are all comparable.\nIn FIG. 1 is shown a typical wafer processing plasma reactor 10. This reactor includes a metal wall 1 1 that encloses a plasma reactor chamber 12. Wall 11 is grounded and functions as one of the plasma electrodes. Gases are supplied to chamber 12 from a gas source 13 and are exhausted by an exhaust system 14 that actively pumps gases out of the reactor to maintain a low pressure suitable for a plasma process. An rf power supply 15 connected to a second (powered) electrode 16 capacitively couples power into a plasma in chamber 12. A wafer 17 is positioned on or near powered electrode 16 for processing. Wafers 17 are transferred into and out of reactor chamber 12 through a port such as slit valve 18.\nA plasma consists of two qualitatively different regions: a quasineutral, equipotential conductive plasma body 19 and a boundary layer 110 called the plasma sheath. The plasma body consists of substantially equal densities of negative and positive charged particles as well as radicals and stable neutral particles. RF power coupled into the reactor chamber couples energy into the free electrons, imparting sufficient energy to many of these electrons that ions can be produced through collisions of these electrons with gas molecules. The plasma sheath is an electron deficient, poorly conductive region in which the gradient in the space potential (i.e., the electric field strength) is large. The plasma sheath forms between the plasma body and any interface such as the walls and electrodes of the plasma reactor chamber.\nWhen the powered electrode is capacitively coupled to the rf power source, a negative dc component V.sub.dc, of the voltage at this electrode (i.e., the dc bias) results (see, for example, H. S. Butler and G. S. Kino, Physics of Fluids, 6, p. 1348 (1963). This dc component, typically several hundred volts, is a consequence of the unequal electron and ion mobilities and the inequality of the sheath capacitances at the electrode and wall surfaces.\nThe dc component of the sheath potential at the powered electrode is useful in accelerating ions to higher energy in a direction substantially perpendicular to the powered electrode. Therefore, a wafer 17, to be processed by ions from the plasma, is positioned on or slightly above the powered electrode 16 so that this flux of positive ions is incident substantially perpendicular to the plane of the wafer.\nThe process rate is dependant on the dc component of the sheath potential and on the density of ions in the plasma near the sheath. To achieve high process throughput, it is advantageous to have an increased density of ions near this sheath. One method of increasing plasma ion density near the wafer utilizes magnets to produce a magnetic containment field that traps electrons within the vicinity of the wafer, thereby increasing the ion production rate and associated density at the wafer. The magnetic containment field confines energetic electrons by forcing them to spiral along helical orbits about the magnetic field lines, caused by a Lorentz force F=q(E+vxB), where E is the electric field vector and B is the magnetic field vector. Unfortunately, nonuniformities of the magnetic containment field of such \"magnetically enhanced\" plasma processing systems result in increased nonuniformity over the surface of the wafer.\nTo reduce the spatial variation of process rate uniformity over the surface of the wafer, some systems, such as the model 5000E from Applied Materials, Inc., produce a magnetic field that is substantially parallel to the surface of the wafer along a direction that rotates around an axis that is concentric with and perpendicular to the wafer. This rotation produces at the wafer surface an approximately cylindrically symmetric time-averaged field that has improved average uniformity over the wafer, thereby producing improved process uniformity.\nAlthough the rotating magnetic field is to produce a uniform processing rate across the entire surface of the wafer, tests of the process rate at various points of the wafer exhibit some significant nonuniformity. It is therefore advantageous to identify the source of this nonuniformity and to provide a mechanism for eliminating or compensating for this nonuniformity."}
{"text": "This invention relates to a biometric authentication system, and more particularly, to registration of a biometric certificate generated from biometric information.\nA biometric authentication system performs personal authentication based on biometric information such as a fingerprint, a vein, an iris, a face, voice, and handwriting. Those pieces of biometric information are private information to be handled with care, and it is therefore required to manage the information properly to prevent leakage of the information. Further, the biometric authentication is expected to be widely used as social infrastructure from now on, but in order to use a large number of biometric authentication systems, it is required to register biometric information with each system. Thus, the user's time and effort to register the information is a factor that inhibits the widespread use of the biometric authentication systems.\nJP 2013-123142 A is known as a technology for solving those problems. JP 2013-123142 A has the following description: “At the time of registration, a biometric signature system embeds a predetermined secret key into a feature quantity of biometric information on a user, and issues a biometric certificate containing a set of the feature quantity and a corresponding public key. At the time of signature, the biometric signature system newly generates a pair of a temporary secret key and a temporary public key for a signature feature quantity of biometric information on the user, creates a signature for a message through use of the temporary secret key, creates a commitment by embedding the temporary secret key into the signature feature quantity, and sets a set of the temporary public key, the signature, and the commitment as a biometric signature. At the time of verification of the biometric signature, the biometric signature system verifies the signature based on the temporary public key, and generates a difference secret key and a difference public key from the biometric certificate, the commitment, and the temporary public key to verify their correspondence.”"}
{"text": "1. Field of the Invention\nThe invention relates generally to user interface applications for autonomous driving systems. More specifically, user interfaces for displaying the status of the autonomous driving system are provided.\n2. Description of Related Art\nAutonomous vehicles use various computing systems to transport passengers from one location to another. Some autonomous vehicles may require some initial input or continuous input from an operator, such as a pilot, driver, or passenger. Other systems, for example autopilot systems, may be used only when the system has been engaged, thus the operator may switch from a manual to an autonomous mode where the vehicle drives itself. These systems may be highly complicated and generally do not provide for a user friendly experience."}
{"text": "A common technique for transporting call control signals over a digital telecommunication network, such as, but not limited to a T-1 data rate network, is the use of `robbed`-bit signaling, in which least significant bits are `robbed` from selected DS0 communication channels and used instead for the transport of in-band signaling information (termed ABCD bits). Although in-band, robbed bit signaling allows the service provider to use all of the available TDM channels for customer traffic, and has been found to be generally acceptable for the transport of quantized voice, it can constitute a significant impairment to the quality of transported data traffic.\nThis data traffic degradation problem can become particularly exacerbated in networks containing a plurality of concatenated signaling discontinuities, such as repeaters and/or cross-connect nodes--that do not allow DS1 extended superframe alignment to maintained. Since the robbed bits cannot be tracked across these discontinuities, a series of three dB noise penalties may be incurred. One possible alternative of using one of the (twenty-four) communication channels as an out-of-band channel for the transport of signaling information for the remaining twenty-three `clear` DS0 communication channels is unacceptable to telecommunication service providers as a cost prohibitive allocation of resources and usurping of useful DS0 bandwidth."}
{"text": "A shallow groove isolation (SGI) structure is now available to make an electrical insulation or isolation between adjacent elements such as transistors, etc. on a semiconductor substrate. As shown in FIGS. 1A to 1D, the SGI structure typically comprises a shallow groove formed on a semiconductor substrate 31 of silicon and an oxide film 35 and the like embedded in the groove and is suitable for devices requiring processing dimensional precision of 0.25 μm or under, because its processing dimensional precision is higher than that of the structure so far by local oxidation of silicon (LOCOS). However, the SGI structure sometimes suffers from formation of sharp protrusions 34 of semiconductor substrate 31 of silicon formed in the oxide film 35 formed by oxidation at the upper edge of the groove during the oxidation step, as shown in FIG. 1C. The presence of such sharp protrusions 34 of semiconductor substrate 31 of silicon causes concentration of electric fields around the protrusions during the circuit operation, sometimes deteriorating gate breakdown voltage or capacitance, as disclosed, for example, by A. Bryant et al (Technical Digest of IEDM '94, pp. 671-674). It is known from experiences that such deterioration of gate breakdown voltage occurs when the radius of curvature of the substrate is not more than 3 nm around the groove upper edge, even if the angle of substrate is not less than 90° around the groove upper edge. To overcome the deterioration, pad oxide film 32 of FIG. 1B is recessed backwards by about 0.1 μm as shown in FIG. 1B′ and oxidized with an oxidant, preferably steam at a temperature of about 1,000° C. to form a desired radius of curvature at the groove upper edges, as disclosed in JP-A-2-260660.\nEven though the desired radius curvature can be obtained by the prior art procedure, step (or unevenness) 44 is formed on the upper surface of semiconductor substrate of silicon 31 around the groove upper edge, as shown in FIG. 1C′. Such step 44 can be formed presumably due to the following mechanism. That is, semiconductor substrate 31 of silicon has a silicon-exposed region and a silicon-unexposed region in the recessed area at the edge of pad oxide film 32; the silicon-exposed region undergoes faster oxidant diffusion, i.e. faster oxidation, than the silicon-unexposed region, resulting in formation of step 44 at the edge of pad oxide film 32 as a boundary. Gate oxide film 37, when formed in such a step region, has an uneven thickness, which leads to variations of electrical properties. Furthermore, stresses are liable to concentrate therein, resulting in a possible decrease in the electrical reliability of a transistors to be formed on step 44.\nFurther, the silicon oxide film 36 is deposited on the semiconductor substrate 31 by chemical vapor deposition (CVD) to embed the silicon oxide film 36 in the groove and then the semiconductor substrate 31 is heat treated to sinter the silicon oxide film 36 embedded in the groove. Sintering is carried out for improving the quality of the silicon oxide film 36 embedded in the groove. If the sintering is insufficient, voids are often generated in the groove in the subsequent steps.\nFurthermore, it is said that wet or steam oxidation is effective for sintering the silicon oxide film 36 embedded in the groove, but the wet or steam oxidation is liable to oxidize the inside, particularly side wall, of the groove. Oxidation starts from the groove surface and thus the groove bottom is less oxidized. Once the groove side wall is oxidized, the active region is narrowed. This is another problem. Thicker oxide film will cause a larger stress on the boundary between the oxide film and the substrate and the once rounded shoulder edge will return to the original sharp one and crystal defects are also generated. This is a further problem. To overcome these problems, it was proposed to provide a silicon nitride film along the groove inside wall.\nAccording to a process for forming a groove, disclosed in JP-A-8-97277, a groove is trenched on a silicon substrate at first, and then an oxide film is formed on the groove inside surfaces (side wall and bottom surfaces) by heat oxidation, followed by further formation of a silicon nitride film thereon and still further formation of a silicon film such as anyone of amorphous, polycrystalline and monocrystalline silicon films on the silicon nitride film. Then, the groove is embedded with a silicon oxide film completely, followed by flattening of the groove top. After the deposition of the silicon oxide film on the entire surface of substrate, but before the fattening, the silicon film is oxidized in an oxidizing atmosphere including steam at about 950° C. to convert it to a silicon oxide film. The silicon substrate is not oxidized during the oxidation, because the silicon substrate is protected by the silicon nitride film. According to the process, a film having a good compatibility with a silicon oxide film, i.e. a silicon film is formed as a thin film on the groove inside surfaces and thus the groove can be embedded with the silicon oxide film without any remaining voids in the groove. The silicon film in the groove must be then converted to a silicon oxide film by oxidation, but the silicon nitride film is provided between the silicon film and the silicon substrate, the silicon substrate is never oxidized during the oxidation of the silicon film. That is, no device characteristics are deteriorated at all.\nIn the above-mentioned prior art processes for forming a groove, heat treatment is carried out at a high temperature such as 1,000° C. or higher to round the shoulder edge of element isolation groove. However, large-dimension wafers are liable to undergo dislocation, which will serve as nuclei for defects, by heat treatment at a high temperature such as 1,000° C. or higher, and thus a heat-treatment process at a high temperature such as 1,000° C. or higher would be hard to use in view of the future trend to use much larger-dimension wafers. In the heat treatment at a low temperature such as less than 1,000° C., it is hard to round the shoulder edge of element isolation groove."}
{"text": "1. Field of the Invention\nThis invention relates to the injection of idle current into a current switch and, more particularly, to the use of fast response diodes to cause the injection of the idle current in the current switch to reduce the propagation delay in said switch.\n2. Description of the Prior Art\nIn the prior art, both Emitter-Coupled Logic (ECL) gate circuits and cascode circuits have been used at current switches. FIG. 1 illustrates the basic current-switch emitter-follower ECL gate circuit. The propagation time delay between the midpoint of a signal V.sub.IN, input to the basic circuit and the midpoint of (either of) its output voltage signal(s) may be expressed as:\nt.sub.pd = t.sub.pd1 + t.sub.pd2 + t.sub.pd3 wherein (for a positive input signal transition), PA1 t.sub.pd1 is the time required to charge the emitter-base junction (capacitance) of the input transistor from its initial value to a value at which the transistor can begin to conduct (e.g. from 0.4 to about 0.7 volts), PA1 t.sub.pd2 is the time required for current transfer given that both switching transistors are fully activated, PA1 t.sub.pd3 is the time required to charge (or discharge) the equivalent capacitance at either collector node through the load resistor connected thereto.\nTo improve the speed of such basic ECL circuits, the possibility of incorporating idle current injection features and \"keep alive\" (KA) diodes in a manner similar to that suggested in FIG. 2 and discussed below was studied by Rigby with the intent of reduction of the t.sub.pd1 component of delay and addition instead of another smaller delay component, the diode switching time. Rigby's work, reported in 1963 in an article entitled High Speed Emitter Current Switching, published in the (Australian) technical journal, Proceedings of IREE, dealt only with the use of discrete component circuits and in general found that the response time of conventional p-n junction diodes available in 1963 and considered by him was too great to afford any meaningful overall circuit delay reduction. One aspect of the invention at hand, however, concerns the use of Schottky or other fast response KA diodes in the basic current switch, their incorporation in integrated circuits being accomplished in a manner similar to that used in the manufacture of \"Schottky Transistor Logic\" so as to entail the expense of relatively little extra monolithic circuit silicon area, and utilizing only currently \"standard\" manufacturing processes. The switching time of such Schottky diodes, being much shorter than t.sub.pd1 in a circuit not incorporating them, meaningful delay reduction can actually be attained.\nAnother aspect of the invention pertains to the use of Schottky or any other KA diodes and idle current injection circuitry in cascode switching circuits, in a manner suggested in FIG. 3 and discussed below. In this application, idle current injection permits improvements of cascode circuit performance in a manner not described in the literature or any known prior art. Neither cascode circuits or their improvements were discussed by Rigby. Such cascode circuits have greater logical flexibility and computational capability than basic ECL gates because of their series gating structure. One of the problems though with these switches is their slightly increased delay in switching relative to conventional current switch emitter follower circuits and their generation of spurious output signals or glytches under certain conditions, as described later below. The present invention though nearly eliminates these glytches and speeds up the switching time of, both lower and upper switching sections, said circuits by injecting idle current through the lower current switch transistors of said circuits and to the cascode nodes, i.e. those at the emitters of the upper switching transistors. The injection of said idle current is controlled by the use of Schottky diodes.\nIn still another aspect of the present invention, no KA diodes are used and idle currents are injected only at the upper current switch cascode nodes. This injection and the means for its accomplishment as well as the particular improvements obtained thereby, not described in any known literature or prior art, are also disclosed in the following."}
{"text": "Some vehicles are used in mines. One such vehicle is generally known as a mine scoop. This versitile vehicle has a low profile for mobility within the low-overhead mine tunnels beyond the area generally known as the cross-cut.\nThe mine scoop has a hydraulic system and as most vehicles, has areas of exposed grease, oil and other moist substances.\nCoal dust is prevalent in the atmosphere of the mine shafts. This highly abrasive, fine dust adheres to moist areas on the mine scoop such as around hydraulic lines where hydraulic fluid may have leaked. The coal dust not only adheres to these moist areas but packs and builds a substantial deposit. As a result, expensive damage may occur to the hydraulic hose on the vehicle unless these coal dust deposits are periodically removed.\nThe existance of the highly explosive atmosphere of the mine beyond the cross-cut requires strong safety measures. For example, cleaning of equipment in this area is usually only done with water. To clean effectively, the scoop must be brought to a less hazardous area where strong cleaning solutions may be used in connection with high pressure cleaning equipment. Portable cleaning equipment is available but, since it is electrically powered, it cannot be used in the areas having a highly explosive atmosphere.\nContamination problems attributed to coal dust are not limited to vehicles but also may do damage to tools and other equipment used in the mine. Cleaning of all such equipment is required.\nIn view of the above, it would be advantageous to provide a washing device using water and detergent under high pressure for cleaning vehicles and other equipment which overcomes the problems associated with the prior art."}
{"text": "1. Field of the Invention\nThis invention relates generally to apparatus for balancing the rotor of a magnetic suspension system including a magnetic bearing and more particularly to apparatus for auto-balancing the rotor so that it spins about its principal axis of inertia rather than its central geometric axis when the two axes are not coincident.\n2. Description of the Prior Art\nIn any suspension system for a rotating element such as the rotor of an electromagnetic machine, rotor balance poses a problem since it is virtually impossible to machine and mount a rotor in bearings so that the axis of inertia of the rotor exactly coincides with the axis of rotation defined by the bearings. The resulting non-coincidence results in undesired vibration and power loss. To alleviate this problem, great efforts have been expended to mechanically balance the rotor with high precision and great delicacy; however, it is virtually impossible to thereafter compensate for aging or thermal deformations without additional mechanical rebalancing requiring undesired shut-down of the equipment and furthermore such deformations cannot always be substantially compensated for all operating speeds.\nWith the development of magnetic bearings for the suspension of a rotor, the existence of any rotor unbalance results in the tendency of the rotor to rotate about the principal axis of inertia lying closest to a desired axis of rotation defined by the bearing rather than the axis of rotation. In the context of this application, the principal axis of inertia of the rotor is hereinafter referred to simply as the axis of inertia or inertial axis. Where such a condition exists, it can be compensated for by detecting the position of the rotor for any departure from its predetermined axial position from which energizing signals are generated and applied to the windings of the bearing to bring the axis of rotation back into proper alignment.\nOne known technique for compensating for synchronous disturbances of a rotor which is supported by a radial magnetic bearing is disclosed in U.S. Pat. No. 4,121,143, entitled, \"Device For Compensating Synchronous Disturbances In The Magnetic Suspension Of A Rotor\", issued to H. Habermann, et al. on Oct. 17, 1978. As disclosed in this patent, a two axis tracking notch filter implemented by a pair of lowpass integrators and two resolvers are connected in a two axis feedback loop which is coupled into the X and Y axis position control loops of the magnetic bearing control system. The tracking notch filter reduces the control loop gain at the rotor's frequency of rotation, thus allowing the rotor to spin about its inertial axis rather than its central geometric axis. This in turn reduces the reaction forces and vibration coupled to the stator and so reduces the power dissipated in the control system. An inherent limitation exists in such a system due to the fact that since resolvers multiply their respective inputs by the factors sin .omega.t and cos .omega.t, where .omega. is proportional to rotational speed of the rotor and generate therefrom a pair of output signals utilized in the feedback loop, the loop gain reduction can only occur at the fundamental rotational frequency .omega.. Accordingly, any imbalances and asymmetries which exhibit higher harmonics of the rotational fundamental frequency, will still transmit vibrations forces to the stator and thus cause the control system to necessarily dissipate power in trying to oppose these forces.\nIt is a primary object of the present invention, therefore, to provide an improvement in the control of a magnetic suspension system for a rotor which is supported by a radial electromagnetic bearing.\nIt is another object of the invention to provide an improvement in the auto-balancing a rotor in a magnetic bearing system.\nAnd yet another object of the invention is to provide an improvement in a system for auto-balancing the magnetically suspended rotor so as to allow the rotor to spin about its inertial axis rather than its geometric axis.\nAnd still a further object of the invention is to provide an improvement in an auto-balancing magnetic bearing system which operates to eliminate external vibration and reduce power consumption."}
{"text": "1. Field of the Invention\nThis invention relates to an image formation apparatus such as a copier or a printer and in particular to an improvement in an image formation apparatus of the type comprising a charger having a charging member in contact with or brought close to the top of a photoreceptor.\n2. Description of the Related Art\nIn recent years, demands for miniaturizing a color image formation apparatus and devises thereof have been made as the demands of the market.\nFor example, a tandem image formation apparatus comprises a plurality of photoreceptors such as photoconductor drums and devices such as a charger and a developing device disposed on each of the photoreceptors. As the apparatus itself is miniaturized, a so-called cleanerless system wherein a cleaning device and a remaining toner collection device are not provided for each of photoreceptors is often adopted.\nHowever, in this kind of cleanerless system, after transfer, remaining toner exists on a surface of the photoreceptor although a trace quantity of toner remains, and the remaining toner becomes xe2x80x9cmemoryxe2x80x9d at a next image formation time, to affect the image quality adversely.\nThus, in a related art, for example, an art has been proposed wherein a memory removal member (for example, a brush roll) is placed in an upstream of a charger member (for example, a charge brush) as a charger to disturb the remaining toner (for example, refer to JP-A-Hei.4-371975).\nHowever, in this kind of image formation apparatus, for example, using a charger of a charging roll type, as shown in FIG. 11, a phenomenon in which random spots are produced at arbitrary points on paper or continuous points, for example, which are spots in correspondence with every rotation period of a charging roll or a photoconductor drum (P/R), are produced is observed. The spots are roughly classified into background spots (BKG spots) occurring in a background and image part spots occurring in an image part (for example, a halftone image).\nParticularly, it is acknowledged that such a spot phenomenon appears noticeably at the initial use stage of first using an image formation apparatus.\nNext, the production principle of such spots is estimated. For example, as shown in FIG. 12, when a foreign substance 502 such as an aggregate of a developer is deposited on a photoconductor drum 510 and enters a nip area between a charging roll 511 and the photoconductor drum 510, the foreign substance 502 portion shields an electric field and a tenting part is formed on a surface layer film portion of the charging roll 511 in which the foreign substance 502 intervenes to cause a charge failure to occur in the part corresponding to the photoconductor drum 510 portion.\nAt this time, if the charge failure part caused by the foreign substance 502 moves to downstream of the photoconductor drum 510 and an electrostatic latent image is formed in the charge failure part and is developed, a spot having a comparatively large diameter is produced.\nOn the other hand, if the foreign substance 502 is deposited, for example, on the charging roll 511 side, a continuous point is produced every rotation period of the charging roll 511.\nPursuing the production cause of such a spot phenomenon, the inventors have estimated that the main cause of producing a spot is the fact that since the developer in a developing device often coagulates particularly at the initial use stage of an image formation apparatus, the aggregate of the developer is easily deposited on a charging roll.\nIf the developing device is immediately used at the initial use stage of the image formation apparatus, shortage of the charge amount of the developer is prone to occur and accordingly foreign substances of an external additive, a coating agent, etc., of the developer are parted to be easily transferred to the photoreceptor side; it is estimated that this factor also becomes the cause of producing a spot.\nThe invention is intended for solving the above-described technical problems and it is an object of the invention to provide an image formation apparatus intended for effectively avoiding a spot phenomenon at the initial use time.\nAccording to the invention, there is provided an image formation apparatus comprising: a photoreceptor; a charger having a charging member placed in contact with or close to the photoreceptor, the charger for charging the photoreceptor; a latent image write unit for writing an electrostatic latent image onto the photoreceptor charged by the charger; a developing device having an agitating and charging element of a developer, the developing device for rendering visible the electrostatic latent image written onto the photoreceptor with the developer; and a performance maintaining controller for executing a performance maintaining initial sequence for agitating and charging the developer uniformly to such a degree that aggregating of the developer in the developing device is at least eliminated under a condition that the image formation apparatus is first used, wherein in the performance maintaining initial sequence, a unit sequence is executed a predetermined number of times, the unit sequence comprising a charge-up mode for agitating and charging the developer in the developing device to raise the charge amount of the developer and a cleaning cycle for cleaning a foreign material of the developer transferred to the photoreceptor.\nAccording to the invention, as shown in FIG. 1, there is provided an image formation apparatus comprising a photoreceptor 1, a charger 2 having a charging member 2a placed in contact with or close to the photoreceptor 1, the charger 2 for charging the photoreceptor 1, a latent image write unit 3 for writing an electrostatic latent image onto the photoreceptor 1 charged by the charger 2, a developing device 4 having an agitating and charging element 4a of a developer G to render visible the electrostatic latent image written onto the photoreceptor 1 with the developer G, and a performance maintaining controller 5 for executing a performance maintaining initial sequence A for agitating and charging the developer G uniformly to such a degree that aggregating of the developer G in the developing device 4 is at least eliminated under a condition that the image formation apparatus is first used.\nIn such technical means, the invention includes not only a form in which a single photoreceptor 1 is provided, of course, but also a tandem form in which a plurality of photoreceptors 1 are arranged and further includes various forms such as a form in which a record material transport body and an intermediate transfer body are placed to face the photoreceptor 1.\nThe charger 2 is required to have the charging member 2a placed in contact with or close to the photoreceptor 1.\nSince it is taken into consideration that it is possible to charge by minute space discharge even in the form in which the charging member 2a is placed close to the photoreceptor 1, the form in which the charging member 2a is placed not to contact with the photoreceptor 1 is also included.\nHowever, preferably the charging member 2a is placed in contact with the photoreceptor 1 because it is facilitated to position the charging member 2a relative to the photoreceptor 1 and the dimension accuracy of the charging member 2a need not be high.\nFurther, the charger 2 may basically comprise the charging member 2a, but is not necessarily be limited to this. For example, the charger 2 maybe provided with a removal member disposed in contact with the photoreceptor 1 in upstream of the charging member 2a, the removal member for removing the deposit on the photoreceptor 1.\nThe removal member may be of contact type for removing the deposit on the photoreceptor 1 or may temporarily remove the deposit on the photoreceptor 1; the removal member may serve as a functional member for eliminating a situation in which the deposit on the photoreceptor 1 arrives at the charging member 2a and holding a good charge property.\nThe developing device 4 is mainly intended for a developing device using a dual-component developer, but is not necessarily be limited to this.\nIn this case, the developer G, which is a spherical toner having a form factor of 130 or less, may be used from the viewpoint of easily providing high image quality and a cleanerless system.\nFurther, the agitating and charging element 4a of the developing device 4 may be any element for agitating and charging a developer and may be an agitating member such as so called an auger, a charge bias with an AC component superposed (in a developing area, the developer alternation-operates to be agitated and charged), etc.\nAn appropriate apparatus may be selected as the performance maintaining controller 5 if the apparatus can execute the performance maintaining initial sequence A.\nThe expression xe2x80x9csuch a degree that aggregating of the developer G in the developing device 4 is at least eliminatedxe2x80x9d means that aggregating of the developer G, which is the main cause of the spot phenomenon at the initial use time, may be eliminated by executing the performance maintaining initial sequence A.\nThe expression xe2x80x9cat leastxe2x80x9d is used to assume that ejecting (corresponding to a cleaning cycle) of a foreign substance of the developer, such as an external additive or a coating agent, which are another main cause of a spot phenomenon at the initial use time, or the like is also contained.\nFurther, the requirement xe2x80x9cagitating and charging the developer G uniformlyxe2x80x9d is based on the fact that not only charging, but also agitating is required to eliminate aggregating of the developer G.\nA representative form of the performance maintaining initial sequence A executed by the performance maintaining controller 5 is to execute a predetermined number of times a unit sequence comprising a charge-up mode for agitating and charging the developer G in the developing device 4 to raise the charge amount of the developer G and a cleaning cycle for cleaning a foreign material of the developer G transferred to the photoreceptor 1.\nIn this case, mainly, the xe2x80x9ccharge-up modexe2x80x9d corresponds to agitating and charging the developer and the xe2x80x9ccleaning cyclexe2x80x9d corresponds to the ejecting of the foreign substances of the developer such as an external additive.\nParticularly, in the form in which the developing bias with an AC component superposed is applied to the developing device 4 in an image formation mode, preferably the performance maintaining controller 5 applies at least a charge bias with an AC component superposed to the developing device 4 in the charge-up mode.\nIn this case, charging up the developer G is promoted by the charge bias with an AC component superposed and ejecting foreign substances in the developer G is also promoted.\nThe charge bias may be any charge bias so long as superposed an AC component and may be set the same as or different from the developing bias.\nIf the user opens interlock by mistake, for example, while the performance maintaining initial sequence A is executed, and the apparatus itself shuts down accordingly, a situation in which the performance maintaining initial sequence A is interrupted can occur.\nAs a preferred form in such a circumstance, the performance maintaining initial sequence A executed by the performance maintaining controller 5 may make it possible to store an incomplete sequence process if the performance maintaining initial sequence A is interrupted at a midpoint and restart the performance maintaining initial sequence A from the midpoint.\nAccording to the form, fruitless repeating the performance maintaining initial sequence A can be avoided.\nThe performance maintaining controller 5 may be any as long as it executes the performance maintaining initial sequence A; it may also execute a performance maintaining sequence B in addition to the performance maintaining initial sequence A.\nFor example, in addition to the performance maintaining initial sequence A, the performance maintaining controller 5 may execute the performance maintaining sequence B for uniformly agitating and charging the developer G in the developing device 4 at least based on elapsed time since the immediately preceding image formation mode.\nThe performance maintaining sequence B is executed to improve a phenomenon in which the charge performance of the developer G is degraded in correspondence with leaving time, and to avoid degradation of the image quality.\nIn this case, as the performance maintaining sequence B, it is common practice to the unit sequence including a charge-up mode and a cleaning cycle a predetermined number of times, but the performance maintaining sequence B is not limited to this.\nFurther, a measuring unit for measuring xe2x80x9cthe elapsed time since the immediately preceding image formation modexe2x80x9d may be a unit for measuring based on the power on time, but preferably the measuring unit can continuously measure the elapsed time when the power is off from the viewpoint of more precisely measuring the elapsed time.\nIn the form, preferably the performance maintaining initial sequence A involves the charge-up mode degree distributed larger than that of the performance maintaining sequence B.\nAccording to this form, the spot phenomenon can be avoided more reliably because the charge property of the developer G can be provided more reliably at the initial use time.\nAt this time, the expression xe2x80x9cthe charge-up mode degree distributed largerxe2x80x9d includes prolonging the charging time itself, enhancing the charging degree itself of the developer G, or a combination both.\nA preferred form of the performance maintaining sequence B may be taken an environmental condition into consideration.\nIn this case, the performance maintaining sequence B may be to uniformly agitate and charge the developer G in the developing device 4 based on the elapsed time since the immediately preceding image formation mode and the environmental condition of the image formation apparatus.\nThe expression xe2x80x9cthe environmental condition of the image formation apparatusxe2x80x9d may be a humidity condition largely affecting the performance of the developer G, but is not limited to this. An ambient temperature condition, temperature and humidity conditions, etc., maybe selected whenever necessary.\nAs a representative determination technique as to xe2x80x9cthe condition that the image formation apparatus is first used,xe2x80x9d for example, in a form in which a process cartridge that can be attached to and detached from the main unit of the image formation apparatus is installed, an identifier indicating that the process cartridge is unused may be provided and the performance maintaining controller 5 may determine whether or not the image formation apparatus is first used based on the identifier information of the process cartridge to execute the performance maintaining initial sequence A.\nThe expression xe2x80x9cthe process cartridge is unusedxe2x80x9d means not only a new product, but also a product unused in a wide sense including a recycled product.\nOn the other hand, the identifier means an unused flag, etc., previously written into a monitor in the process cartridge, for example."}
{"text": "It is well known that ambient illumination, that is light originating from sources external to the display device, is reflected to the observer from various optical interfaces of the device and thus reduces the image contrast by increasing the apparent brightness of the dark image areas. Under conditions of high ambient illumination, the image contrast is severly degraded. In addition, a part of the light emitted by the luminescent material of the device also undergoes undesired reflections, producing a further degradation of contrast and of resolution. When the luminescent material consists of a layer of phosphor material in the form of small powder particles, scattering of the emitted light also occurs, further degrading resolution.\nVarious means for overcoming these problems have been proposed. These include the use of various filters including polarizing, neutral density and restricted angle or multi-apertured opaque filters. Other methods include the incorporation of a dark material into the glass of the tube face, or a black dye in the phosphor dielectric layer of the display device. All of the methods have the common disadvantage that the emitted light as well as the reflected ambient light intensity is reduced, with the result that the improvement is contrast ratio is less than desired because the emitted light intensity is a factor upon which the contrast ratio depends.\nThe remarkable reflection-reducing properties of inhomogeneous films were recognized as early as 1880 by Lord Rayleigh (Proc. Lond. Math. Soc. 11, 51, 1880); the properties of such films have been extensively reviewed in a recent series of articles by Jacobsson (Progr. in Optics 5, 247, 1965; Arkiv Fysik 31, 191, 1966; Physics of Thin Films 8, 51, 1975). According to Jacobsson, experimental studies to date have been mainly devoted to transparent inhomogeneous films composed of graded mixtures of two nonabsorbing materials such as ZnS--Na.sub.3 AlF.sub.6, ZnS--CeF.sub.3, CeO.sub.2 --CeF.sub.3, and CeO.sub.2 --MgF.sub.2. These films were found to be durable and of good optical quality. A high index mixture of Ge--ZnS has been produced for application in the infrared wavelength region but were found to be relatively soft and sensitive to moisture and inferior to Ge--MgF.sub.2 films. KBr--Au films were found to have a very low absorption index, with k = 0.01 even at a concentration of gold of 0.16 parts by volume of gold. By contrast, an absorption index of 1.0 was found for a Ge--Au mixture containing 0.1 parts by volume of gold. Ge--In films were also found to have relatively high absorption. Due to the low solubility of In in Ge, the In was expected to remain a separate phase in the form of more or less spherical inclusions.\nAn inhomogeneous Ge--Si.sub.x O.sub.y film was shown by Jacobsson (1965) and also Olsen and Brown (Res./Develop. 16, 52, 1965) to lower the reflectance of a Ge surface to that of a surface of Si.sub.x O.sub.y (refractive index 1.62). Even lower reflectance was obtained with Ge--MgF.sub.2 films, although the transmittance was higher than expected (Jacobsson and Martensson, App. Optics, 5, 29, 1966). One of the first applications of inhomogeneous films as an antireflection coating was described by Nadeau and Hilburn in Canadian Pat. No. 418,289 (1944), and U.S. Pat. No. 2,331,716 (Oct. 12, 1944), in which a plastic layer of polystyrene or urea-formaldehyde resin having a high refractive index is diffused into the surface of an article and overcoated with a second plastic of low refractive index such as cellulose caproate or ethylcellulose. An important commercial application of inhomogeneous films as a low reflectance, absorbing coating on sunglasses was described by Anders in U.S. Pat. No. 3,042,542 (German Pat. No. 1,075,808; 1960). The inhomogeneous films described by Anders consisted of a mixture of low refractive index material, CeF.sub.4, ThF.sub.4, MgF.sub.2, or SiO.sub.2, and a metal, Ni, Fe, Mn, or Cr, or lower oxide of Nb, Ta, or Ti.\nRecently, Steele has proposed in U.S. Pat. No. 3,560,784 the use of a dark dielectric layer consisting of SiO.sub.2 with a tapered concentration of codeposited aluminum applied to the rear side of a light transmissive phosphor layer to serve as a light absorbing layer. The tapered concentration of aluminum results in a continuous variation of the index of refraction through the layer, and such layer comprises an optically inhomogeneous film. Steele claimed novelty for a high contrast cathode ray tube utilizing this construction in which the refractive index of the silicon oxide was substantially equal to that of the phosphor. Phospors suitable for use with the inhomogeneous film of Steele were not otherwise identified. The same objective was the object of an earlier patent of Coltman (U.S. Pat. No. 2,616,057) in which the light absorbing layer was described as lampblack or the black deposits produced by evaporating metals such as aluminum or antimony under poor vacuum conditions.\nUp to the present, the deposition of tapered inhomogeneous films such as in the Steele patent has required the evaporation of two different materials, with the rate of evaporation of each varied as a function of time. Also, it is usually desired that the initial portion of the deposit consist of one component only with the end portion consisting of the second different material only. Steele shows the initial and end materials to be SiO.sub.2 and aluminum, respectively. These requirements pose severe technical difficulties and to achieve reproducible results, elaborate monitoring and control equipment is required so that despite the superior performance offered by inhomogeneous films as compared to homogeneous films, very limited commerical application has been made of inhomogeneous films.\nOsterberg (J. Opt. Soc. Am. 48, 513, 1958) has shown that transmitted waves cannot suffer loss of energy by reflection as they traverse nonabsorbing, inhomogeneous media in which the optical properites have no discontinuities. This result is strictly true only when the medium is infinite in extent. For practical applications, film thicknesses used are of the order of the wavelength of light so that interference due to reflection at the boundaries occurs. The width of the reflectance minimum has been found, however, to be greater than can be achieved with homogeneous films. It also has been shown by Osterberg that inhomogeneous absorbing media similarly cannot exhibit reflectance when the optical properties are continuous. In this case, the medium need not be infinite in extent. Anders (Dunne Schichten fur die Optik, Wissenschafftliche Verlagsgesellschaft mbH, Stuttgart, 1965, English translation as Thin Films in Optics, The Focal Press, London, 1967) has observed that a film thickness of only one wavelength is sufficient for essentially complete absorption in an absorbing inhomogeneous film. This property is basic to the dark dielectric layer described by Steele in U.S. Pat. No. 3,560,784 (1971) since the tapered concentration of aluminum results in an absorbing inhomogeneous film. The deposition of such film entails, however, the technical difficulties previously described, including the deposition of two different materials from two sources."}
{"text": "1. Field of the Invention\nThis invention relates to non-volatile semiconductor memory capable of storing an analog value or multiple bits per memory cell.\n2. Description of Related Art\nA typical electrically erasable non-volatile memory contains a memory array including hundreds or thousands of rows of memory cells and hundreds or thousands of columns of memory cells, where each memory cell contains a transistor such as a floating gate or split-gate transistor having a programmable threshold voltage. The threshold voltage of a memory cell (or a transistor within the memory cell) indicates a stored value and is programmed by applying appropriate programming voltages to the control gate, source, and drain of the memory cell. Depending on the resolution of read and write circuits which read and write threshold voltages, non-volatile memory can store one bit, an analog value, or a multi-bit value per memory cell.\nIn a typical non-volatile memory architecture, each row of memory cells in a memory array has a row (or word) line coupled to control gates of the memory cells in the row, and each column of memory cells has a column (or bit) line coupled to drains of the memory cells in the column. Sources of the memory cells in an array may be connected in a variety of ways. For example, a virtual ground array often has neighboring columns of memory cells that share a column line so that each column line in a virtual ground array is coupled to the drains of memory cells in one column and to the sources of memory cells in the neighboring column. For flash memory, sources of all memory cells in an erasable sector are coupled to a source line for the sector. Accessing memory cells for erase, write, or read operations applies appropriate voltages to the row, column, and source lines coupled to the memory cell or cells to be erased, written, or read. When programming (or writing to) a selected memory cell, programming voltages applied to selected row, column, and source lines for a sufficient period change the threshold voltage of the memory cell coupled to those lines.\nA problem in non-volatile memories is that voltages applied to row and column lines of memory cells selected for a write or erase operation can disturb or change the threshold voltages of memory cells in the same columns or rows as the selected memory cells. The accumulated threshold voltage disturbances from writing to hundreds or thousands of memory cells in the same row or column can significantly change the threshold voltage of a memory cell and the value stored in the memory cell. This problem is particularly significant for large memory arrays which have more memory cells per row and column. The problem is also significant in analog or multilevel memories where a relatively small variation in a threshold voltage changes the value stored. Accordingly, methods and architectures for reducing the write and erase disturb in analog and multilevel memories are sought."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of soybean breeding. In particular, the invention relates to the novel soybean variety D5852641.\n2. Description of Related Art\nThere are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include higher seed yield, resistance to diseases and insects, better stems and roots, tolerance to drought and heat, better agronomic quality, resistance to herbicides, and improvements in compositional traits.\nSoybean, Glycine max (L), is a valuable field crop. Thus, a continuing goal of plant breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties."}
{"text": "Anxiety and depression are major psychiatric disorders of significant clinical and socioeconomic significance. Clinical Depression generally presents alongside Anxiety Disorders, and vise-versa. Rarely does a patient present symptoms of only one or the other.\nIn the general population, these disorders affect daily performance and correlate with impulse control, financial behaviors, substance abuse and organization. Anxiety is an unpleasant state that involves a complex combination of emotions that include fear, apprehension, and worry. It is often accompanied by physical sensations such as heart palpitations, nausea, chest pain, shortness of breath, or tension headache. Anxiety disorder is a blanket term covering several different forms of abnormal, pathological anxiety, fears, phobias and nervous conditions that may come on suddenly (acute anxiety) and/or gradually over a period of several years (chronic), and may impair or prevent the pursuing of normal daily routines. Anxiety disorders are often debilitating chronic conditions, which can be present from an early age or begin suddenly after a triggering event. They are prone to flare up at times of high stress.\nAnxiety is often described as having cognitive, somatic, emotional, and behavioral components (Seligman, Walker & Rosenhan, 2001). The cognitive component entails expectation of a diffuse and uncertain danger. Somatically the body prepares the organism to deal with threat (known as an emergency reaction): blood pressure and heart rate are increased, sweating is increased, bloodflow to the major muscle groups is increased, and immune and digestive system functions are inhibited. Externally, somatic signs of anxiety may include pale skin, sweating, trembling, and pupillary dilation. Emotionally, anxiety causes a sense of dread or panic and physically causes nausea, and chills. Behaviorally, both voluntary and involuntary behaviors may arise directed at escaping or avoiding the source of anxiety. These behaviors are frequent and often maladaptive, being most extreme in anxiety disorders. However, anxiety is not always pathological or maladaptive: it is a common emotion along with fear, anger, sadness, and happiness, and it has a very important function in relation to survival.\nNeural circuitry involving the amygdala and hippocampus is thought to underlie anxiety (Rosen & Schulkin, Psychol. Rev., 105(2):325-350, 1998). When confronted with unpleasant and potentially harmful stimuli such as foul odors or tastes, PET-scans show increased bloodflow in the amygdala (Zald & Pardo, PNAS, 94(8):4119-4124, 1997; Zald, Hagen & Pardo, J. Neurophysiol., 87(2):1068-1075, 2002). In these studies, the participants also reported moderate anxiety. This might indicate that anxiety is a protective mechanism designed to prevent the organism from engaging in potentially harmful behaviors.\nConventional treatments for anxiety include behavioral therapy, lifestyle changes and/or pharmaceutical therapy (medications). Most drugs used to treat these disorders are known to have negative side effects that may limit their use, or cause habituation and dependence.\nPostsynaptic density-95 protein (PSD-95) couples NMDARs to pathways mediating excitotoxicity and ischemic brain damage (Aarts et al., Science 298, 846-850 (2002)). This coupling was disrupted by transducing neurons with peptides that bind to modular domains on either side of the PSD-95/NMDAR interaction complex. This treatment attenuated downstream NMDAR signaling without blocking NMDAR activity, protected cultured cortical neurons from excitotoxic insults and reduced cerebral infarction volume in rats subjected to transient focal cerebral ischemia. This result has led to the proposal to use peptide antagonists of PSD-95/NMDAR for treating stroke and other diseases mediated by excitotoxicity. No significant side effects have been observed in phase I trials of one such antagonist."}
{"text": "Many wireless network operators enter into roaming agreements with one or more other wireless network operators. The roaming agreements allow wireless devices to exchange wireless communications with visited wireless networks when a home wireless network, to which the wireless devices are subscribed, is not available. Moreover, a roaming agreement may allow a visited wireless network to provide wireless connectivity over a protocol that is not available on a home wireless network in the same area despite the home wireless network offering connectivity in the area over other protocols. For example, the visited wireless network may provide service in the area using both 3G and 4G wireless protocols while the home wireless network only provides service using the 3G protocol. Thus, even though the home wireless network is available in the area, a wireless device may be able to roam on the visited wireless network using the 4G protocol.\nThe roaming agreements discussed above typically require a fee be paid the home wireless network operator to the visited wireless network operator for the use of the visited wireless network. Moreover, these fees may vary depending on the protocol used by a wireless device when roaming. Accordingly, limiting when a wireless device roams and what protocol the wireless device uses when roaming may help minimize roaming fees paid by the home wireless network operator.\nOverview\nEmbodiments disclosed herein provide systems and methods for compelling a cell selection metric to indicate that a wireless cell is not suitable for communications. In a particular embodiment, a method provides receiving system information from a base station of a visited wireless communication network, wherein the system information includes at least a portion of cell selection parameters for a wireless cell of the base station. The method further provides calculating a cell selection metric from the cell selection parameters and a cell selection offset stored in the wireless communication device, wherein the cell selection offset causes the cell selection metric to indicate that the wireless cell is not suitable for communications. The method further provides receiving a request for a data connection from an application executing on the wireless communication device and, in response to the request, recalculating the cell selection metric from the cell selection parameters without the cell selection offset. The method further provides determining whether the wireless cell is suitable for communications based on the cell selection metric."}
{"text": "1. Field\nThe invention pertains to enhancing the quality of recorded service data, such as data recorded on service tickets, in a data center or call center.\n2. Description of the Related Art\nService delivery centers are large, complex and dynamic ecosystems, which engage hundreds of thousands of experts globally to manage thousands of processes supporting thousands of IT systems with hundreds of configurations. While operations at service delivery centers are typically associated with back-end processes, its efficiency affects quality at front-end (e.g., client experience and satisfaction).\nMultiple ticketing systems, data stores and warehouses trace the operations in service delivery centers. They capture practices of Subject Matter Experts (SMEs), who are typically System Administrators (SAs), and changes in the IT infrastructure (e.g. server decommissioning). These ticketing systems, and enterprise-level warehouses are only reliable as their sources, whether human-driven (tickets submitted by SAs) or system-driven (automated updates of server registries).\nAll too often, there is poor quality of captured data when managing a data center or call center. Administrators are time pressured to achieve high throughput and problem resolution, and no incentive exists for quality of records and logs when capturing and describing problems and resolutions. Low quality of such data leads to inefficiencies in operations (e.g. incomplete tickets slow down the problem resolution process), or leads business analytics to reach wrong or suboptimal conclusions. Frequently, data records such as tickets are blank with insufficient data, and as such are unusable.\nMoreover, low quality of data affects the business decisions (e.g. leading to poor business insights when identifying opportunities for new service offerings, such as “show me the low utilization servers across the banking sector”). Business insights and problem resolution processes require careful quality assessment to build credibility with stakeholders and efficiently resolve problem tickets. Moreover in such volatile environments, quality of operations and business insights will vary depending on the corresponding data source.\nPlanning activities also depend on good quality data. Take for example server consolidation, where old servers or underutilized servers are migrated into virtual environments with newer hardware. Being able to understand the configuration information such as number of CPUs, speed, memory, operating system and software configured as well as resource information such as network bandwidth, disk and CPU utilization are all key to be able to prepare a plan that maps to proper sized servers. Bad quality data could easily derail a plan from improper source selection to bad target allocations.\nAccumulated problem resolution records contain tremendous source of information about the managed system, its efficiencies and weaknesses, and in addition to analytics, it is a valuable source for knowledge transfer and learning in attempt to train new administrators. The record data are also used for reporting and report generation in billing and service level agreement (SLA) measurements.\nAccurate records of services provided are valuable for a number of business aspects. These include planning of future system improvements, automating problem resolution, optimization of tasks, and awarding the best administrators and skill development. It would be desirable to have a way to improve capturing of incident and problem description and resolution in a data or call center."}
{"text": "With the advent of Internet Protocol TV (IPTV), consumers may receive and manage media programming functions such as video and/or music on demand from an IPTV set-top box (STB) receiver, computer, digital telephone system, or other computing device. The IPTV STB receiver can also be used for digital video recording (DVR) and to management personal media files such as picture albums and family movies. IPTV devices, however, are not readily manageable when the consumer is in transit.\nA need therefore arises for a system and method for configuring media services."}
{"text": "Panoramic videos enable a viewer to view video content in an immersive video environment.\nThe listing or discussion of a prior-published document or any background in this specification should not necessarily be taken as an acknowledgement that the document or background is part of the state of the art or is common general knowledge. One or more aspects/examples of the present disclosure may or may not address one or more of the background issues."}
{"text": "1. Field of the Invention\nThe invention relates to an arrangement for microwave transmission between wave guide regions having different internal gas pressures and/or different fill-gas compositions, that is to say, for coupling or outcoupling microwaves of such a wave guide region into another region.\n2. Description of the Related Art\nIn German Patent Application No. DE-OS 36 22 614 which corresponds to commonly owned U.S. Pat. No. 4,877,642, is disclosed a method of manufacturing electrically conductive moulded bodies by a plasma-activated chemical deposition from a gaseous phase. With such methods the coupling of high-power microwaves is effected through a hermetically sealed insulating microwave aperture of dielectric material in a microwave resonator used as a reaction chamber, in which a plasma is formed and electrically conductive layers are chemically deposited. During this process the problem arises that an electrically conductive film generally covers the surface of the microwave aperture arranged at the coupling place, that is, its inside surface facing the reaction chamber, as a result of which the coupling is stopped. This problem is solved according to DE-OS 36 22 614 either by having the inside of the microwave aperture rinsed by an inertial gas, or selecting for the microwave aperture a dielectric material which is kept free from growth of electrically conductive film as a result of an etching reaction with one of its reaction partners.\nA cognate problem occurs when high-power microwaves from gyrotrons are outcoupled during transition from high-vacuum to air. With microwave powers of the order of 0.1 to 1 MW the thermal load of the known materials used for microwave apertures becomes too large, as a result of which the output power is restricted. With maximum power levels of 0.3 MW one manages by enlarging the wave guide and additionally cooling the aperture consisting of, for example, Al.sub.2 O.sub.3.\nEvacuation of a wave guide through non-radiating or non-coupling slots is known from British Patent Specification No. GB-PS 644,749."}
{"text": "The rapid rise of health care costs has become an important issue in modern society. To help reduce the costs, professional care givers have begun to seek alternatives, one of which is home health care services. These services not only tend to reduce costs, but are also preferred by the patient wishing to remain in his familiar environment. Among the many types of services provided are: respiratory care, rehabilitation therapy, cardiac monitoring procedures, and infusion therapy.\nInfusion therapy involves IV administration of drugs. Making this therapy safe and convenient for a home situation allows a great number of patients who would otherwise be hospitalized to remain at home and still receive medication. Currently, over 300,000 patients annually use a home infusion therapy delivery system. Typically, patients include the elderly with chronic diseases like cancer, patients with either Crohns disease, HIV or other immune system disorders, and patients suffering from chronic pain. Many of these patients require infusion treatment over a long duration such as months or even years.\nOne characteristic of home IV drug therapy, in contrast to hospital administered therapy, is that a nurse is not always present or readily available. To provide safe and effective treatment, home infusion therapy usually requires that the patient himself, or other non-professional caregiver, such as a relative, administer IV fluids. Special training is required because many home care patients on IV therapy require multiple drugs or multiple doses of the same drug each day. The average nursing visit to a home infusion therapy patient is typically about 90 minutes including commuting time. The typical patient gets between 1 and 4 nursing visits per week, but has to take IV medications daily. Since the cost of daily care by a nurse is not usually covered by most insurers, the cost of attention by a nurse is most economically applied in training the patent or other amateur caregiver and in monitoring the therapy program.\nIn the home care situation non-compliance, over-medication or under-compliance with the IV therapy protocol is a serious issue and quite prevalent. For instance, non-compliance (not taking a medication) or under compliance (taking fewer or smaller dosages than prescribed) occurs in up to approximately one-third to one-half of elderly home therapy patients. Typical compliance related problems include forgetting to follow the specified procedure for administration of the IV medication, forgetting to turn on the various devices used to administer the IV medication and forgetting to turn off a medical device which then delivers too much medication (over-medication). Reasons for compliance related problems are varied and include poor communication, confusion or forgetfulness regarding the procedures and/or equipment, or even attempts to avoid the adverse side effects of IV medications and fluids. Misapplication of the home IV therapy protocol can have serious ramifications resulting in greatly increased home health care nursing expenses, re-hospitalization, and reduction in health status of the patient. Thus, there is a strong need for improved monitoring of patient compliance with the health care program. Benefits of such improved monitoring and compliance include, but are not limited to, improved health at a lower cost, while still remaining in the preferred home environment.\nTo properly monitor compliance with an IV therapy protocol, a device may be provided for monitoring the flow of IV medications and fluids. The IV fluids for a single patient are likely to come from several different sources or systems including IV pumps, IV fluid controllers, gravity drips, syringes, and other devices.\nA typical gravity powered IV may be as simple as an IV bag hanging on a pole in which a nurse or care giver manually adjusts a valve to limit the flow rate, but not control it accurately, or it may use an electronic controller which optically counts the drops of fluid as they pass an optical sensor and then adjusts the flow rate accordingly. However, optical drop counting sensors only provides an indication that the fluid is flowing past the sensor when in a vertical orientation such as hanging from an IV pole. Thus the patient and IV delivery equipment must remain relatively stationary during the administration of the medication or fluid. Optical drop counters also function poorly at higher flow rates and higher line pressures, such as when a syringe is used, because the fluid moving past the drop counter tends to become a continuous stream rather than remaining discrete drops. Therefore, the optical drop counter technique cannot be adapted for use with all fluid sources.\nAn alternative to an optical drop counting sensor, or as a stand-alone measuring device, is a single point pressure transducer to measure the fluid pressure in the IV tubing at a selected point of measurement. This type of sensor is common in IV pumps and is used to indicate that the pump is generating a static pressure head and, correspondingly, causing fluid flow or backpressure in the event of an occlusion in the IV line. This type of sensor only determines line pressure at the selected point, and is only useful in monitoring the pressure caused by the IV pumping device and the related backpressure caused by moving fluids into the patient's body. However, this type of single-point pressure sensor is useful in many IV delivery systems to determine if fluid pressures are at correct levels, and to detect changes in fluid pressure which are indicative of an occluded or collapsed vein. Often, when a certain threshold pressure is detected in a device using this type of sensor, an alarm is sounded to warn of a flow problem. This type of device measures changes in the static line pressure of a fluid line, but is unable to determine if a patient is following proper IV drug administration procedures and cannot differentiate between changes in pressure due to fluid flow versus some other cause, such as an occlusion in which there is actually no fluid flow.\nIncreased backpressure in an IV fluid line causes problems, and, as described above, many IV fluid delivery systems use a sensor to determine when high backpressure develops, i.e. , greater than a few inches of water, for instance when an infiltration of tissue occurs or the tubing becomes occluded. Upon the detection of a significant backpressure, the device sounds an alarm and may function to automatically discontinue the delivery of the IV medication and fluids. Therefore, it is important that any device used to monitor whether or not fluid is flowing does not cause a substantial increase in backpressure or a false occlusion alarm might be triggered.\nOther alternatives use indirect methods to monitor the flow of IV fluids. For instance, the speed and number of rotations in a pump mechanism may be monitored to indirectly determine when fluid flow is occuring. This is useful for flows caused by an IV pump, but is of no value to patients who also receive gravity drips or fluids via syringe. Since nearly all infusion therapy patients must perform venous access device maintenance procedures, such as a heparin flush via syringe to maintain the patency of their IV lines, this pump rotation technique is not of value for monitoring all infusions.\nThe time usage for an IV delivery system may be recorded to prepare bills to patients. Typically, the information is printed or stored in an electronic memory device such as the electronic controls for the drop counter or IV pump. The information may be used to determine which of several patients are using the IV system being monitored, it may be used to coordinate several IV delivery systems with a centrally managed pump, or it may be used to facilitate billing and reimbursement. Unfortunately, none of these systems accommodate tracking of fluid delivered from a variety of sources such as to a patient who receives syringes, gravity drips, and IV pump infusions. The present invention provides an improved flow indicator switch, which overcomes the above-mentioned limitations of the prior art."}
{"text": "1. Field of the Invention\nThis invention is related to a camera for close-up photography using a photo-ranging system to insure that the subject to be photographed is properly focused. The ranging system is used in combination with a selected one of a plurality of fixed focus lenses mounted at the front of the camera.\n2. Description of the Prior Art\nCameras having the capacity to take close-up photographs are not new. An example of a camera used for close-up photography is the Acmel Macro Auto V6 camera manufactured by Acmel Corporation of Tokyo, Japan which involves a highly complicated system involving focal beams and a plurality of detachable lenses for different focal lengths. Each detachable lens is designed for a single subject distance. The resulting camera serves its purpose but is altogether too complicated for easy use by a camera operator. It is both bulky and heavy, weighing about 41/2 pounds.\nThe use of converging light beams from a camera to determine a proper subject to camera distance is old art as disclosed in U.S. Pat. No. 3,416,426 which is specifically incorporated herein by reference. The theory disclosed therein and in other literature provides a light source reflected from spaced apart mirrors through a lens board to converge at the subject of the photograph. This technique is useful where several parameters remain constant, in particular, one of the constants is flash photography. With flash photography the duration of exposure is set and the only other variables are the exposure aperture and the focus of the lens.\nU.S. Pat. No. 4,777,501 discloses a single pair of laser ranging lights of different wavelengths mounted to project converging light beams at the focal point of a single fixed focus exposure lens of the associated camera.\nU.S. Pat. No. 4,914,460 discloses a single pair of light sources each of which projects a plurality of light beams in a particular pattern. The pattern is displayed on a surface which allows an observer to determine the topography of the surface and its distance from an associated camera. The camera is mounted equidistant between the two light sources.\nU.S. Pat. No. 5,142,299 discloses an underwater camera mounted between a pair of light sources. The light sources project light beams to converge at the focal point of a single fixed focus lens system of the associated camera.\nWhat is missing from the prior art is a camera for close-up photography having more that one focal length available to a user combined with a plurality of ranging systems to insure proper focus of the subject."}
{"text": "Within the past four years, the genetic cause of the Hereditary Nonpolyposis Colorectal Cancer Syndrome (HNPCC), also known as Lynch syndrome II, has been ascertained for the majority of kindreds affected with the disease (13). The molecular basis of HNPCC involves genetic instability resulting from defective mismatch repair (MMR). To date, six genes have been identified in humans that encode for proteins and appear to participate in the MMR process, including the mutS homologs GTBP, hMSH2, and hMSH3 and the mutL homologs hMLH1, hPMS1, and hPMS2 (2, 7, 11, 17, 20, 21, 22, 24). Germline mutations in four of these genes (hMSH2, hMLH1, hPMS1, and hPMS2) have been identified in HNPCC kindreds (2, 11, 13, 17, 24). Though the mutator defect that arises from the MMR deficiency can affect any DNA sequence, microsatellite sequences are particularly sensitive to MMR abnormalities (14). Microsatellite instability is therefore a useful indicator of defective MMR. In addition to its occurrence in virtually all tumors arising in HNPCC patients, Microsatellite instability is found in a small fraction of sporadic tumors with distinctive molecular and phenotypic properties (27).\nHNPCC is inherited in an autosomal dominant fashion, so that the normal cells of affected family members contain one mutant allele of the relevant MMR gene (inherited from an affected parent) and one wild-type allele (inherited from the unaffected parent). During the early stages of tumor development, however, the wild-type allele is inactivated through a somatic mutation, leaving the cell with no functional MMR gene and resulting in a profound defect in MMR activity. Because a somatic mutation in addition to a germ-line mutation is required to generate defective MMR in the tumor cells, this mechanism is generally referred to as one involving “two hits,” analogous to the biallelic inactivation of tumor suppressor genes that initiate other hereditary cancers (11, 13, 25). In line with this two-hit mechanism, the non-neoplastic cells of HNPCC patients generally retain near normal levels of MMR activity due to the presence of the wild-type allele."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device and a method of manufacturing the same.\n2. Description of the Related Art\nA micro-electro-mechanical system (MEMS) display is a display expected to replace a liquid crystal display (see Japanese Patent Application Laid-open No. 2008-197668). This display differs from a liquid crystal shutter type display utilizing polarization, and displays an image by opening and closing a light transmissive window using a mechanical shutter system. A shutter is formed of a thin film. Vertical and horizontal sizes of one shutter forming one pixel are in the order of several hundred micrometers, and a thickness thereof is in the order of several micrometers. One shutter is opened/closed to enable ON/OFF operation for one pixel. The shutter is operated by an electrostatic attractive force.\nThe shutter is arranged in a space surrounded by a sealing member between a pair of light transmissive substrates, and the space is filled with oil. The oil is used to prevent a spring for driving the shutter from sticking, and to reduce a difference in optical refraction index with respect to the light transmissive substrates.\nThe oil is injected from an injection port corresponding to an opening of the sealing member, and the injection port is encapsulated by a resin after the oil injection. For high-speed open/close operation of the shutter, oil having low viscosity is desired, but in this case, the speed of oil when passing through the injection port increases to cause a problem of damaging the shutter."}
{"text": "One or more aspects relate, in general, to processing within a processing environment, and in particular, to optimizing the processing.\nProcessors execute instructions that direct the processors to perform specific operations. The instructions may be part of user applications that perform user-defined tasks, or part of operating system applications that perform system level services, as examples.\nOne processing technique used by the processors to process the instructions is referred to as pipelined processing, in which processing is performed in stages. Example stages include a fetch stage in which the processor fetches an instruction from memory; a decode stage in which the fetched instruction is decoded; an execute stage in which the decoded instruction is executed; and a complete stage in which execution of the instruction is completed, including updating architectural state relating to the processing. Other and/or different stages are possible.\nTo facilitate processing within a pipelined processor, various optimization techniques are employed. One such technique includes decode time instruction optimization, which offers an opportunity to improve code execution by combining multiple instructions into a single internal instruction; recombining multiple instructions into multiple/fewer internal instructions; and/or recombining multiple instructions into multiple internal instructions with fewer data dependencies."}
{"text": "The present invention relates to a computer-implemented method for signing a message by a user device of a public key infrastructure (PKI) system. The present invention further relates to a corresponding attestation server and a corresponding computer program product.\nThe de-facto technique for authenticating messages are digital signatures. Digital signatures allow the holder of a private key to generate a signature which can be verified using the corresponding public key. Such digital signatures are based on the property that no one except the holder of the private key can generate signatures that are valid under the public key. Digital signatures may be used e.g. during TLS client authentication and server authentication, for signing contracts and for e-mails.\nTo verify that a given public key belongs to a given entity, public key certificates may be used. Such certificates effectively bind a public key to a given entity. A certificate may be generated by a trusted authority of a PKI-system. The trusted authority, which is often denoted as certificate authority, issues a signature on the public key of the respective entity and on additional attributes corresponding to this entity, e.g., the name, the address and/or its affiliation. This signature, together with the information to verify it, then acts as the certificate.\nIn order for the digital certificate system to be secure, the private key of the signer must be kept secret. If the private key gets compromised, e.g., if a computer has been infected by malware, or a smart phone is lost, the certificate on the corresponding public key has to be revoked. This means that the certificate issued beforehand is not valid anymore and a verifier will reject a signature signed under the corresponding public key. The revocation may be done by informing the issuer of the certificate (certificate authority) and requesting a revocation of the certificate. The certificate authority may publish certificate revocation lists.\nTo verify a signature, a verifier must determine the validity of both the certificate and the signature itself. Hence usually a public key certificate is sent along with the corresponding digital signature. The verifier may determine whether a certificate is still valid by asking the certificate authority whether the certificate has been revoked and/or by checking revocation lists that have been published by the certificate authority. The validity of the certificate should be verified at each verification of the given signature.\nAccordingly, there is a need for alternative methods for verifying signatures in a PKI-system."}
{"text": "Information is easily dispersed through the Internet, television and many other outlets. One major problem is that the information dispersed is often not correct. Although there are fact checking websites available online, these websites check facts in a slow manner; typically not truly providing a fact check response for several hours or even days."}
{"text": "1. Field of the Invention\nThe present invention relates to a rocker-type electrical switch suitable for commercial and home use.\n2. Description of the Prior Art\nKnown is a rocker-type electrical wall switch which comprises a rocker pivotally supported in a housing at a first pivot point, a movable contact brush pivotally supported at a second pivot point in the housing, a spring compressed between a downwardly extending boss on the rocker and a lower end of the contact brush, the spring being movable under compression to inclined positions relative to the brush in response to pivotal movement of the rocker between rest positions, the movement of the spring transmitting pivotal movement of the rocker to the brush, and a pair of spaced cams engaging, respectively, with the upper end of the brush at a point above the second pivot point, and the rocker and cams being movable into engagement with the brush under pressure exerted by the spring on the rocker.\nOther known devices of some relevance to the present invention are one which discloses a safety snap switch; one which teaches a snap switch based on the engagement between a rigid member, able to oscillate, and a resilient prestressed contact in such a manner that rebound is substantially prevented; one which teaches a number of toggle type switches having various contact structures; one which teaches a switch including a contact-carrying rocker, the movement of which is produced by a compression spring, the axis of which coincides with that of a control knob or a lever, the spring transmitting its action to the rocker through a link or stirrup engaging through its end on the one hand, the rocker, and on the other hand, the spring; one which teaches a snap-action electrical switch with contact dampening means to quiet the action of lever-operated electric switches; one which teaches an electrical toggle switch having a mounting that can oscillate for the contact in the inner position and association of the mounting with a simple form of an essentially leaf-type spring; one which teaches a noiseless electric switch having a pivoted operating lever biased into two switch positions by a leaf spring and the lever; and one which teaches a compact electrical contact and electrical switch structure having a combination of a screw terminal, a push-in wire terminal, and a make or break electrical contact terminal, with the three terminals being formed in a single compact electrical structure from a small piece of metal strip bent at right angles between the screw terminal and the push-in terminal.\nAlso known is a device which comprises a mounting strap for supporting a wiring device in a metal wall box and establishing an electrical connection between the metal mounting screw and the strap. The mounting screw is inserted through the strap and threaded into a metal box or gem box."}
{"text": "1. Field of the Invention\nThe present invention relates to means for transferring an ignition impetus through the wall of a firearm to a confined propellant, and more particularly to a nipple for use in conjunction with firearms.\n2. Description of the Prior Art\nSince the early use of firearms, it has been necessary to transfer an ignition impetus from the exterior of the barrel or cylinder of a firearm to a propellant charge, such as black powder, disposed therein. Early in the 19th century, flintlocks, which touched off a small black powder charge were employed and this was later largely replaced by the use of percussion caps which, upon impact, produce a desired quantity of burning gas or \"flash\" which is employed to ignite the propellant charge. These percussion caps are customarily fitted to what is known as a nipple which is treated through an aperture disposed into the barrel of a rifle or hand gun, adjacent to the breech portion thereof, or in the rearward wall of a cylinder in a revolver. When the percussion cap is snugly in place on the nipple, the force of striking the percussion cap not only produces a flash but also forces this gas under considerable pressure through the nipple into the ignition chamber of the firearm.\nTypically, the interior portion of the nipple terminates in an aperture. Unfortunately, a frequently incurred problem is the blockage of this aperture for various reasons. The result is the placing of a percussion cap on the nipple and the striking thereof with the end result of no ignition of the propellant charge.\nVarious methods of clearing such a blockage have developed over the years and there is keen interest in this problem today as the use of antique or reproduction of antique black powder firearms has become a large hobby. Solutions include the removal of the nipple and the clearing of the aperture in which the nipple was disposed, the forcing of a fine wire through the nipple and the somewhat, dangerous practice of removing the nipple, placing a small amount of fresh black powder adjacent to the location of the interior opening of the nipple, in the propellant, and replacing of the nipple. Unfortunately, none of these methods is totally satisfactory and, if a powder charge is in place covered by a lead ball or bullet it has been known to be necessary to employ a screw-type extractor or the like to remove the ball and permit removal of the powder charge so that a barrel or cylinder can be totally cleared for reuse with fresh powder.\nIn an attempt to solve other problems associated with these early firearms, various configurations were proposed for extending the flashhole in the nipple through the powder charge. U.S. Pat. No. 36,464 issued to Hopkins on Sept. 16, 1862; U.S. Pat. No. 21,802 issued to Schenkl on Oct. 12, 1858; and U.S. Pat. No. 15,292 to Halsey on July 8, 1856 teach the employment of a tige or tube disposed in the barrel of a rifle so that the point at which the flash from a percussion cap or the like touches off the powder is extended to a point somewhere remote from the point of entry of the flash into the barrel. These configurations are provided so that the powder can be ignited entirely before it is blown out of the rifle barrel. Further, in Halsey it is suggested that the tige principle can be used in conjunction with a nipple wherein the tige is a mere continuation of the nipple what is deemed a proper distance into the interior of the barrel to form a tube. Although these references teach methods of enhancing combustion, the basic problem of an orifice which can be clogged by a powder charge still exists.\nThe present invention overcomes the problems associated with the prior art by providing a nipple which provides an extended opening from which an ignition impetus such as that generated by a percussion cap can contact the powder thereby limiting the possibility of the extended orifice becoming clogged.\nOther percussion nipples are shown in U.S. Pat. No. 4,163,335 to Ives and U.S. Pat. No. 4,186,506 to Pawlak. In addition, a percussion nipple is also shown in German Auslegeschrift No. 1,216,018 to Hintze. However, none of the references cited show or suggest the configuration of the present invention which has solved an extremely long standing problem."}
{"text": "For the purpose of prenatal genetic diagnosis, amniotic diagnosis has heretofore been conducted primarily by sampling amniotic fluid via amniocentesis and inspecting chromosomes of the fetal cells in the amniotic fluid. Conventional prenatal genetic diagnostic techniques suffered from the serious problems of a risk of miscarriage in addition to physical and mental stresses on mothers. Under such circumstances, fetal cells (fetal nucleated red blood cells) were found to migrate in the blood circulating in the mother's body. If fetal nucleated red blood cells contained in the maternal blood are selectively collected and the genes of the fetus are analyzed, prenatal diagnosis can be safely carried out without a risk of miscarriage. Such technique enables fetal gene diagnosis at an early stage of pregnancy, which can lead to early treatment. Approximately 5,000,000 cases of prenatal genetic diagnosis are conducted every year on a global scale. If such safe genetic diagnostic technique can be put to practical use, safe techniques can occupy a high share of the global market. However, it is not easy to collect fetal nucleated red blood cells because such cells are said to exist in amounts as small as about 1 cell in 1 ml of the maternal blood. A collection method involving the use of an antibody that recognizes a special structure of the nucleated red blood cell surface (i.e., an antigen-antibody reaction), a method comprising allowing fluorescence-labeled nucleated red blood cells to flow in a liquid, allowing such blood cells to pass through the laser beam focal point, and assaying fluorescence emitted by blood cells to collect cells (i.e., fluorescence activated cell sorting (FACS)), and other techniques have been implemented in research institutes all over the world. However, all such techniques have been insufficient. As a method for collecting nucleated red blood cells with high assuredness, a method comprising analyzing an image observed under an optical microscope and collecting the nucleated red blood cells detected can be employed. According to the FDD-MB® (Fetal DNA diagnosis from maternal blood) project of Takabayashi (Kanazawa Medical University), at present, fetal nucleated red blood cells are separated from the maternal blood via density-gradient centrifugation using Percoll to prepare samples and automatically processed to collect NRBC (Haruo Takabayashi, Idenshi Igaku (Gene & Medicine), Vol. 5, No. 3, 2001, pp. 10-11; Haruo Takabayashi, Idenshi Igaku (Gene & Medicine), Vol. 5, No. 3, 2001, pp. 28-34 2). Detection of the nucleated red blood cells via imaging disadvantageously necessitates a long period of time.\nRare cells have heretofore been separated by density-gradient centrifugation using Ficoll, Percoll, Polymorphprep, or the like (US Patent Publication No. 2003/0134416; US Patent Publication No. 2004/0142463; U.S. Pat. No. 5,714,325; U.S. Pat. No. 6,949,355; U.S. Pat. No. 7,166,443; WO International Publication No. 2008/048931). When such separation reagent is used alone, disadvantageously, nucleated red blood cells cannot be completely separated because their density (specific gravity) is similar to that of white blood cells and some other red blood cells (i.e., 1.07-1.08)."}
{"text": "Design of semiconductor devices such as power semiconductor devices requires trade-offs between electric characteristics such as area-specific on-state resistance Ron×A, breakdown voltage Vbr between load terminals such as source and drain, switching behaviour and device ruggedness.\nBy way of example, increasing a specific resistance of a bulk material allows to achieve lower electric field strengths at a device front side. Although lower electric field strengths at the device front side may improve device ruggedness, a softness of the switching behaviour may be adversely affected.\nIt is desirable to improve the trade-off between electric characteristics in semiconductor devices."}
{"text": "Aspects of the present disclosure relate generally to wireless communication, and more particularly, to methods and apparatus for power management.\nWireless communication networks are widely deployed to provide various communication services such as voice, video, packet data, messaging, broadcast, etc. These wireless networks may be multiple-access networks capable of supporting multiple users by sharing the available network resources. Examples of such multiple-access networks include Code Division Multiple Access (CDMA) networks, Time Division Multiple Access (TDMA) networks, Frequency Division Multiple Access (FDMA) networks, Orthogonal FDMA (OFDMA) networks, and Single-Carrier FDMA (SC-FDMA) networks.\nA wireless communication network may include a number of eNodeBs that can support communication for a number of user equipments (UEs). A UE may communicate with an eNodeB via the downlink and uplink. The downlink (or forward link) refers to the communication link from the eNodeB to the UE, and the uplink (or reverse link) refers to the communication link from the UE to the eNodeB.\nIn some wireless communication networks, a user equipment (UE) selects and maintains a connection with a base station providing communication capabilities for the UE. Further, small cells may be deployed to improve wireless network communications when experiencing poor macro base station connections. In such wireless communication networks, inefficient utilization of available communication resources, particularly power resources, may lead to degradations in wireless communication. Even more, the foregoing inefficient resource utilization inhibits network devices from achieving higher wireless communication quality. In view of the foregoing, it may be understood that there may be significant problems and shortcomings associated with current power management technology. Thus, improvements in wireless network power management are desired."}
{"text": "Transportation systems such as railways can be complex systems, with several components being interdependent on other components within the system. Attempts have been made in the past to optimize the operation of a particular component or groups of components of the railway system, such as for the locomotive, for a particular operating characteristic such as fuel consumption, which can be a significant component of the cost of operating a railway system. Some estimates indicate that fuel consumption is the second largest railway system operating cost, second only to labor costs.\nFor example, U.S. Pat. No. 6,144,901 proposes optimizing the operation of a train for a number of operating parameters, including fuel consumption. Optimizing the performance of a particular train (which may be only one component of a much larger system that includes the railway network of track, other trains, crews, rail yards, departure points, and destination points), however, may not yield an overall system-wide optimization or improvement of one or more of the operating parameters.\nOne system and method of planning at the railway track network system is disclosed in U.S. Pat. No. 5,794,172. Movement planners such as this are primarily focused on movement of the trains through the network based on business objective functions (BOF) defined by the railroad company, and not necessarily on the basis of improving performance or a particular performance parameter such as fuel consumption. Further, the movement planner may not extend the improvement down to the train (much less the consist or locomotive), nor to the railroad service and maintenance operations that plan for the servicing of the trains or locomotives.\nThus, there does not appear to be recognition that improvement of operations for a transportation system may require a multi-level approach, with the gathering of key data at several levels and communicating data with other levels in the system.\nPowered systems that operate within transportation systems or other systems can include off-highway vehicles, marine diesel powered propulsion plants, stationary diesel powered systems, and rail vehicle systems, e.g., trains. Some of these powered systems may be powered by a power unit, such as a diesel or other fuel-powered unit. With respect to rail vehicle systems, a power unit may be part of at least one locomotive and the rail vehicle system may further include a plurality of rail cars, such as freight cars. More than one locomotive can be provided with the locomotives coupled as a locomotive consist. The locomotives may be complex systems with numerous subsystems, with one or more subsystems being interdependent on other subsystems.\nAn operator may be onboard the powered system (such as a rail vehicle) to ensure proper operation of the powered system. In addition to ensuring proper operation of the rail vehicle, the operator also may be responsible for determining operating speeds of the rail vehicle and in-vehicle forces within the rail vehicle (e.g., forces between coupled powered units such as locomotives and/or non-powered units such as cargo cars or other railcars). To perform this function, the operator may have extensive experience with operating the rail vehicle over a specified terrain. The experience and knowledge of the operator may be needed to comply with prescribed operating speeds that may vary based on the location of the rail vehicle along a route, such as along a track. Moreover, the operator also may be responsible for ensuring in-vehicle forces remain within acceptable limits.\nEven with knowledge to ensure safe operation, the operator may not operate the vehicle so that the fuel consumption, emissions, and/or travel time is reduced or minimized for each trip. For example, other factors such as emission output, environmental conditions like noise or vibration, a weighted combination of fuel consumption and emissions output, and the like may prove difficult for the operator to both safely operate the vehicle while reducing the amount of fuel consumed by the vehicle, reducing the amount of emissions generated by the vehicle, and/or reducing the travel time of the vehicle. The varying sizes, loading, fuel characteristics, emission characteristic, and the like can be different for various vehicles, and external factors such as weather and traffic conditions can frequently vary.\nOwners and/or operators of off-highway vehicles, marine diesel powered propulsion plants, and/or stationary diesel powered systems may realize financial benefits when the powered systems produce increased fuel efficiency, decreased emission output, and/or decreased transit time so as to save on operating costs while reducing emission output and meeting operating constraints, such as but not limited to mission time constraints."}
{"text": "It is known that arthroprostheses used to reconstruct the shoulder joint can be the anatomical or inverse type. Anatomical humeral prostheses are provided with a hemispherical head that artificially reconstructs the human anatomy and are suitable to articulate in the glenoid cavity of the shoulder blade or possibly in a mating artificial seating attached to the glenoid cavity. On the contrary, inverse prostheses provide a concave articular insert able to allow the rotation of an artificial spherical body attached to the glenoid seating, commonly called glenosphere.\nIt is known that humeral prostheses, whether anatomical or inverse, generally include a distal joint element, also called rod or pin, inserted inside the humerus along the diaphyseal axis, to support the humeral head in anatomical humeral prostheses or the concave articular insert in inverse prostheses.\nDistal joint elements of various types are known, for example with an elongated tubular shape, to define an elongated rod, or a small-sized pin, typically used for mini-invasive humeral prostheses. The shape and size of the distal joint element can also depend on the type and size of the humeral prosthesis, for example inverse humeral prostheses can need a distal joint element that is bigger than anatomical ones.\nIt is possible, for example due to degeneration of the tissues of the shoulder, that operations may have to be carried out to revise the prosthesis, for example to replace the prosthesis, or operations to convert an anatomical prosthesis into a inverse prosthesis or vice versa.\nIt is known, for example from the patent application EP-A-1.472.999, to make a modular humeral prosthesis, comprising an adapter body, which functions as a metaphyseal module, to be positioned in a seating made under the head of the humerus and connectable on one side to a rod inserted in the humerus and on the other side to the semi-spherical humeral head, in the case of anatomical prostheses, or to the concave articular insert in the case of inverse prostheses. Indeed, by varying the adapter body, the known modular humeral prosthesis can be made as anatomical or inverse. Moreover, in the field of the same type of prosthesis, anatomical or inverse, it may be possible to change the shape or size of the adapter body, releasing it from the rod and keeping the latter inserted in the humerus, whenever its sizes or characteristics are suitable for the humeral prosthesis to be grafted.\nHowever, one possible disadvantage of this known modular humeral prosthesis is that it cannot be converted from anatomical to inverse without having to replace the adapter body. This latter operation, which entails extracting the adapter body from the bone seating, can be invasive and destructive.\nAnother possible disadvantage of the known humeral prosthesis is that it does not allow to completely replace the rod with another element inserted in the humerus, such as a rod with a different conformation or size, or a nail for osteosynthesis, preserving the adapter body already inserted, that is, without extracting the adapter body from the humerus which, as we said above, can be difficult and invasive.\nIt is known for example that, in the event of fractures of the humerus, for the purposes of osteosynthesis, it may be necessary to introduce inside the humerus a humeral nail that reproduces the correct alignment of the bone fragments to allow reciprocal welding thereof. The humeral nail can be introduced in antegrade fashion, with a passage through the proximal part of the humerus (antegrade nail), or retrograde, through an aperture in the olecranon fossa (retrograde nail). The antegrade solution is generally preferable, since it is less invasive and less complex.\nHowever, in the state of the art and with the humeral prostheses available, if there is a fracture of the humerus and there is a humeral prosthesis present, it may be necessary to completely remove the prosthetic implant, or alternatively to use external synthesis means, such as osteosynthesis plates which in any case do not always guarantee a successful synthesis.\nOne disadvantage of known solutions is therefore that it it is not possible to act in an antegrade manner to remove the rod or pin and insert the nail without removing the adapter body, in the event of fractures of the humerus where there is a humeral prosthesis present.\nDocument US-A-2005/0125067 describes a modular prosthesis of the known type, provided with a head and a rod coupled to the head. The rod has a proximal portion coupled with the head and a distal portion configured to extend in a long bone of the patient. The distal portion can be removed from the rod after the prosthesis has been implanted, without removing the proximal portion.\nDocument WO-A-99/47081 also describes a modular orthopedic prosthesis having a body that has a through hole that receives a connection bushing that in turn receives a rod having a proximal neck and a distal shaft. The connection bushing can be expanded radially in order to clamp the rod and the body together.\nPurpose of the present invention is to obtain a modular humeral prosthesis that allows revisions or conversions of the prosthesis, keeping the adapter body of the humeral prosthesis in the humeral seating and allowing to replace the distal joint element, either rod or pin."}
{"text": "1. Field of the Invention\nThe present invention relates to a hot water supply system of connecting a plurality of water heaters, setting a priority device from the water heaters, and performing hot water supply operation while the priority device links with some of the other water heaters, a water heater therefor and a hot water supply control method therefor.\n2. Description of Related Art\nConventionally, there is an art of connecting a plurality of water heaters and performing the hot water supply operation.\nAn example of performing hot water supply using connected water heaters like the above is that a connecting unit is connected to a hot-water supplier capable of being operated by one set individually; when a control unit of the hot-water supplier detects the connection of the connecting unit, the control unit switches the operation mode of the hot-water suppliers from an individual operation mode into a connecting operation mode automatically (for example, Japanese Laid-open Patent Publication No. 2002-357361).\nAnother example thereof is that two hot water supply equipments are connected via a communication cable, and an anomaly in one hot water supply equipment is notified at the other hot water supply equipment (for example, Japanese Laid-open Patent Publication No. 2007-078327).\nWhen many, for example, at least three water heaters are connected and the hot water supply operation is performed, a connecting unit for controlling the hot water supply operation is separately provided in order that the operation of each water heater is linked with each other. In a hot water supply system like the above, it is possible to obtain so much hot water that a single water heater cannot obtain. However, newly disposing the connecting unit causes an increase in cost. Further, a place for the connecting unit is needed.\nAlso, the number of connectable water heaters is limited when the connecting unit is not provided, water heaters are connected to each other directly, one is treated as a master and the others are treated as slaves for example, and the hot water supply operation using the linked water heaters is controlled.\nConcerning such problems, there is no disclosure or suggestion thereof in Japanese Laid-open Patent Publications Nos. 2002-357361 and 2007-078327, and no disclosure or suggestion about structure etc. for solving them is presented."}
{"text": "1. Field of the Invention\nThe present invention relates to a composite material having a reinforcing texture in refractory fibers and ceramic matrix with multiple interphases between the fibers of the texture and the matrix.\nBy refractory fibers is meant here fibers in carbon or in ceramic such as, for example the fibers constituted essentially of silicon carbide.\n2. Prior Art\nIn the composite materials with refractory fibers and ceramic matrix, the fiber-matrix bond conditions the transfer of load and, as a result, the characteristics as well as the mechanical behavior of the materials.\nOne advantageous way of controlling the load transfer between fiber and matrix consists in interposing a fine layer of a material having a lamellar structure oriented in parallel to the axis of the fiber, as described in U.S. Pat. No. 4,752,503. The layer having the lamellar structure is in a material selected from rough laminar pyrolytic carbon (RL) PyC and boron nitride (BN). As shown diagrammatically in FIG. 1, the layer having a laminar structure constitutes an interphase LI defining two interfaces: one interface between fiber F and the lamellar interphase LI and one interface between said latter and the matrix M.\nWithout the lamellar interphase, a crack starting in the matrix M spreads directly through the fiber F, as shown in FIG. 2, this leading to a premature breaking of the fiber. The material has a fragile behavior.\nWhen, on the contrary, there is a lamellar interphase, this prevents any cracks starting in the matrix from spreading directly through the fiber. Due to its resiliency under shear stress, the lamellar interphase permits a relaxing of the stresses exerted on the bottom of cracks. The material has a non-fragile behavior and improved mechanical properties, as shown by curve II in FIG. 3 which shows the relation between elongation and tensile strength. By way of comparison, curve I in FIG. 3 shows this relation in the case of a material without lamellar interphase.\nThe crack remains stopped in the lamellar interphase as long as the level of stress exerted on the bottom of the crack does not exceed the breaking strength of the weakest of the elements found in the immediate vicinity of the crack. Three elements have to be taken into account: the material constituting the lamellar interphase, the matrix-lamellar interphase interface and the fiber-lamellar interphase interface. Depending on which one of said elements has the lowest breaking strength, the progression of the crack will follow path a, b, or c, respectively, when the stress on the bottom of the crack increases, as illustrated in FIG. 4.\nThe most dangerous progression is that following path c, namely on the fiber-lamellar interphase interface. Indeed, the crack then can reach into the fiber and break it, if it meets with any surface defect of said fiber, which will cause a reduction of the mechanical properties of the composite material."}
{"text": "In telecommunication, teleconferencing is the live exchange and mass articulation of information among persons and machines remote from one another but linked by a telecommunications system, for example, a telephone system. Computers have given new meaning to the term because they allow groups to do much more than just talk. Once a teleconference is established, the group can share applications and mark up a common whiteboard.\nBroadly speak, teleconferencing comprises various ways by which people communicate with one another over some distance. In a narrow sense, a teleconference is a two-way, interactive meeting, between relatively small groups of people (approximately 1 to 10 at each end), who may use permanent teleconferencing facilities. A teleconference involves audio communication between the locations, but may also involve video or graphics. One problem with conventional teleconferencing systems is that as more participants are added to the teleconference, the conventional teleconferencing systems' quality and performance degrades. In other words, as more participants are added to conventional teleconferencing systems, the conventional system's overall latency increases and long delays are created between when participants can speak."}
{"text": "The present invention relates to a process for regulating to desired values the dimensions of bubbles of magnetic bubble elements during their production by liquid phase epitaxy.\nIt is pointed out that a magnetic bubble element is constituted by a magnetic layer with small magnetic domains having an opposite magnetic induction to that of the material surrounding them in the layer.\nIn a monocrystalline magnetic layer, such as a magnetic garnet film having a uniaxial magnetic anisotropy perpendicular to the plane of the layer, it is possible to create generally cylindrical magnetic domains in which the magnetic induction is of the opposite direction to that in the remainder of the layer.\nThese domains, which are normally \"bubbles\" are stabilized at their operating size under the action of a continuous magnetic field, called the polarization field. The latter must be perpendicular to the layer and the domains can be displaced in the plane of the layer under the action of propagation means magnetized by a rotary magnetic field applied in the plane of the layer. In this way, it is possible to produce circuits, comparators, memories, etc.\nMore specifically, the present invention relates to the preparation of magnetic bubble elements constituted by ferrimagnetic garnet films deposited by liquid phase epitaxy on a non-magnetic garnet substrate, said films preferably being magnetized perpendicular to the plane of the film. In such films, the magnetic domains appear in the form of cylinders with a circular cross-section, for example, positive on the upper face of the layer and negative on the lower face from the magnetic standpoint and in this way they form magnetic dipoles having an axis perpendicular to the displacement plane.\nIn connection with the construction of magnetic bubble memories, it is known that their capacity is directly linked with the diameter of the magnetic bubbles. Thus, to obtain a capacity of 256 kbits, elements are used, whose bubbles have a diameter of 2.7 .mu.m, whilst to obtain capacities up to 0.5 and 1 megabit it is necessary to use elements whose magnetic bubbles have a diameter of 3 to 1.5 .mu.m, preferably 2.5 and 1.8 .mu.m.\nThus, in connection with the construction of magnetic bubble memories, considerable importance is attached to processors making it possible to adjust the diameter of the magnetic bubbles of such elements to the desired values.\nHitherto, in the processes for the production of magnetic bubble elements by liquid phase epitaxial deposit, the diameter of the element bubbles has been controlled by acting on the composition of the epitaxy bath comprising oxides or carbonates of the elements used in the composition of the film.\nThus, in the case of garnet films of formula: EQU T.sub.a.sup.1 T.sub.b.sup.2 T.sub.c.sup.3 Ca.sub.d) (Fe.sub.e Ge.sub.f)O.sub.12\nin which T.sup.1, T.sup.2 and T.sup.3, which differ from one another, represent an element in the series of rare earths including yttrium and a, b, c and d are numbers such that their sum is substantially equal to 3, whilst e and f are numbers such that their sum is substantially equal to 5, the diameter of the bubbles has been controlled by modifying the composition of the epitaxy bath with respect to the quantities of the different rare earths oxides for influencing the anisotropy and on the quantity of germanium oxide for modifying the magnetization. (Materials Research Bulletin, Vol. 10, No. 1, 1975 and Journal of Crystal Growth, Vol. 12, No. 1, December 1977).\nThus, in processes for the production of bubble memories as desired in the Journal of Crystal Growth, Vol. 12, No. 1, December 1977, certain conditions must be respected in order to obtain films with a satisfactory quality.\nThus, for reducing the diameter d of bubbles of the element, it is necessary to reduce the characteristic length l of the film on wishing to respect the condition according to which said diameter d is similar to the thickness h of the film in order to obtain a good stability of the bubbles. This can be achieved by increasing the magnetization of the film because the characteristic length l is defined by the formula: ##EQU1## in which A represents the exchange constant, Ku the uniaxial anisotropy constand and M.sub.s the saturation magnetization.\nHowever, on increasing the saturation magnetization M.sub.s of the film, the anisotropy field H.sub.k is generally reduced making it difficult to respect the condition: EQU H.sub.k -4.pi.M.sub.s .gtoreq.700 oersteds\nnecessary for preventing spontaneous nucleation of the bubbles. In addition, in order to respect this condition, it is necessary to influence the respective quantities of the rare earths to increase the anisotropy field H.sub.k.\nHowever, this control method involving on the one hand the respective quantities of the different rare earths and on the other the germanium quantity in the epitaxy bath has the disadvantage of requiring a relatively large change in the epitaxy bath composition to change from one bubble diameter to another."}
{"text": "1. Field of the Invention\nThe present invention relates to a resonator, to a vibrating sensor including such a resonator, and to a method of fabricating the resonator. Such a resonator is used for example in vibrating sensors of the gyro type.\n2. Brief Discussion of the Related Art\nThe general principle of a vibrating sensor is to subject a resonator to vibration and to detect a physical magnitude that is representative of the influence of an acceleration on the vibration.\nVibrating sensors exist that include an electrode-carrying plate with a resonator mounted thereon. The resonator comprises a body having a substantially hemispherical resonant part with a pole that is connected to the electrode-carrying plate by a sensor stem. The resonant part comprises a hemispherical web defined by an outside surface and an inside surface, which surfaces have free edges that are connected to each other by a plane annular surface that extends facing the electrodes secured to the electrode-carrying plate. The resonant part and the stem are covered in an electrical conduction layer. The electrical conduction layer and the electrodes are connected to different potentials so as to subject the web to periodic elliptical deformation and so as to detect the orientation of the ellipse, e.g. as a function of capacitances of values that depend on the gap formed between the electrodes and the plane annular surface.\nThe body is made of silica because of its isotropic properties and its very low mechanical damping.\nIn an embodiment that used to be widespread, the resonant part and the stem were completely covered in a layer of chromium forming the conduction layer (see for example document US-A-2003/0019296).\nNevertheless, it was subsequently found that the layer of chromium contributed non-negligible mechanical damping.\nResonators were thus considered in which the layer of chromium did not extend over all of the resonant part: the layer of chromium then covered the stem and the annular surface, and was provided with branches extending over the inside surface from the stem to the annular surface. That enabled damping to be reduced significantly, thereby giving rise to a large increase in the performance of vibrating sensors incorporating such resonators.\nSubsequently, replacing the chromium in those resonators with platinum has enabled the performance of vibrating sensors to be further improved.\nThe document EP-A-1445580 discloses a resonator having an inside surface fully covered by a conductive metal coating and an outside surface left uncovered.\nNevertheless, silica is a material with a very active surface that tends to establish bonds with its surroundings. Silica becomes covered on its surface in groups of silanol Si—OH or silane Si—H that are strongly polar and that can combine with examples in the surroundings. In vibrating sensors, the resonator is in a vacuum so that such contamination of the silica does not occur. Nevertheless, the research that has led to the invention has shown that such contamination of the silica that is left uncovered by the conduction layer gives rise to unstable and anisotropic modification of the geometrical properties and of the mechanical damping of the resonator, even when the level of contamination is low: such modification is likely to lead to drift that cannot be compensated electronically because it is unstable."}
{"text": "Patent literature 1 discloses a cloud-key-management-type decryption technology that stores, in a key device, registered permission information corresponding to a terminal device to which a decryption authority is given, allows the key device to receive ciphertext and terminal information, and outputs response information corresponding to the decrypted result of the ciphertext when the terminal information corresponds to any of the registered permission information."}
{"text": "The present invention is directed to a method of optimizing the steering assistance of a motorized vehicle, using angle sensors instead of a torque detector. The method of the present invention also provides an improved steering power in case of failure. The present invention also encompasses a vehicle comprising two angle sensors used to optimize the steering assistance.\nFor utility vehicles, a steering assistance is necessary. It is usually provided through a torsion bar, which opens a hydraulic valve, according to the torque applied by the driver to the steering wheel. In case of failure of the hydraulic pump, or another part of the steering system, the effort to steer the steered axle considerably increases. In case such a failure occurs on an heavy truck, the driver becomes unable to steer the steerable wheels. It is therefore necessary to provide a backup steering system, which allows at least partial steering power. A back up steering system usually requires a second torsion bar, which is costly, heavy and space consuming. It is sometime not possible to implement such a second torsion bar on the steering column. DE102004049038 describes the use of two angle sensors to record the data resulting from the torsion of the torsion bar. However, DE102004049038 is not directed to backup steering systems.\nIt is therefore desirable to provide a method of optimizing the steering assistance with a costly efficient and space saving solution.\nThe steering system of an aspect of the present invention comprises one torsion bar and two angle sensors. The first angle sensor is positioned upstream the torsion bar and the second angle sensor is positioned downstream the torsion bar, in such a way that the torsion angle of the torsion bar can be monitored by the means of the two angle sensors. The portion of the steering column which is upstream the torsion bar comprises all the mechanical elements between the steering wheel and the part just above the torsion bar. It encompasses for example the upper shaft, the lower shaft, with inner shaft and outer shaft, a steering wheel adjustment device. The portion of the steering column which is downstream the torsion bar encompasses all the elements between the torsion bar and the steered wheels. This part comprises for example the drop arm, ball joints, drag link, the upper steering arm, the track rod. In case of twin steered axles, the portion which is downstream the torsion bar also encompasses the elements involved in the steering of the second steered axle. In particular, the second steering pump, the steering actuator of the second steered axle, and the secondary steering rod are downstream the torsion bar.\nIn a first embodiment, the angle sensors are used to detect an abnormal increase of angle between the first and the second angle sensor.\nThe method of the present invention comprises the steps of\na) Monitoring the steering angle of the steering wheel, by the means of a first angle sensor;\nb) Monitoring the steering angle of the steered wheels, by the means of a second angle sensor;\nc) Comparing the difference between the steering angle of the steering wheel, monitored is step a), and the steering angle of the steered wheels, monitored in step b), with a first reference value and/or comparing the steering angle of the steered wheels, monitored in step b), with a second reference value;\nd) Detecting whether the difference between the steering angle of the steering wheel monitored in step a) and the steering angle of the steered wheels monitored in step b) reaches the first reference value of step c) and/or whether;\ne) If the difference between the steering angle of the steeling wheel monitored in step a) and the steering angle of the steered wheels monitored in step b) reaches the first reference value of step c) and/or the steering angle of the steered wheels, monitored in step b) differs from the second reference value of step c), then activating a failure mode.\nIn step a), the angle to which the driver steers the steering wheel is determined by the means of the first angle sensor, positioned upstream the torsion bar. Each angle of rotation of the steering wheel may be associated or not associated to a theoretical angle of rotation of the steered wheels. The theoretical angle of rotation of the steered wheel is the angle expected for a given steering angle of the steering wheel. It may be for example a linear function of the steering angle of the steering wheel. Alternatively, the theoretical angle of the steered wheels may be a non-linear function of the angle of rotation of the steering wheel. The first angle sensor is preferably an angle sensor already present on the vehicle and involved in other functions. For example, the first angle sensor may be the angle sensor already used for the ESP functions.\nIn step b), the effective steering angle of the steerable wheels is determined by the means of a second angle sensor, positioned downstream the torsion bar. This second angle sensor is preferably positioned close to the torsion bar, on the output shaft of the steering gear, in order to provide a direct measurement. However, the second angle sensor may be positioned anywhere else downstream the torsion bar. In case of twin steered axles, the second angle sensor is preferably positioned on the first steered axle. The second angle sensor is preferably an angle sensor already present in the vehicle and involved in other functions. Indeed, an angle sensor may already be present for the steering management of the second steered axle. In this case, there is no need for additional specific sensors.\nStep a) is concomitant with step b). This means that the steering angle of the steering wheel, is determined in step a) at the same time the steering angle of the steered wheels is determined in step b). Monitoring the steering angles in steps a) and b), or the difference of angles, has to be understood as repeating the operation of determining the steering angles, either permanently or as soon as one of the steering angles is modified. Permanently determining the steering angles means that a regular measurement is performed, for example at a predetermined frequency. Preferably, the steering angle is determined each few milliseconds, most preferably between 1 and 10 milliseconds.\nIn step c), the difference between the steering angle of the steering wheel and the steering angle of the steered wheels is monitored and compared to a predetermined value, which is a first reference value, or a warning threshold value, under which should remain the difference of steering angles. If a theoretical value is associated to the steering angle of the steering wheel in step a), the effective steering angle of the steered wheels, measured in step b), may also be monitored and compared to this theoretical value, which is a second reference value. Under normal conditions, the effective steering angle of the steered wheels should correspond to the second reference value. Also, under normal conditions, the difference of the steering angles determined in steps a and b) should remain under the first reference value. Under these circumstances, it is considered that the suitable steering assistance is delivered, allowing effective steering of the steered wheels. No additional steering power is triggered.\nIn step d), it is identified that the difference of the steering angles, reaches the first reference value or the effective steering angle of the steering wheels departs from the second reference value. Under these conditions, it is considered that the steering system is in fault and step e) is initiated. Alternatively, step e) may be initiated if the two conditions of step b) are reached. In this case, step e) is initiated only when the difference of the steering angles reaches the first reference value and the effective steering angle of the steering wheels departs from the second reference value.\nStep e) triggers a failure mode, wherein additional power steering is delivered to compensate the efforts of the driver. The failure mode may be the activation of an auxiliary steering power. In case of more than one steered axle, the failure mode may be a special mode of the steering system of the second steered axle. For example, under failure mode, the steering system of the second steered axle may be activated in a way to provide an oversteering of the second steered axle. The failure mode may encompass any other action which aims at improving the steering assistance."}
{"text": "The present invention relates generally to the manufacture of book covers and particularly to the controlled and gradual establishment of contact between the adhesively coated cover cloth and cover boards during assembly of a book cover. More specifically, this invention relates to improved methods of and apparatus for fabricating book covers.\nBook covers consist of a covering material, front and back cover boards and a spine board positioned between the two cover boards. The boards are glued to the covering material over their entire area. The covering material that projects beyond the cover profile defined by the boards is turned in on all four sides to define flaps. In the industrial manufacture of book covers, these manufacturing operations are accomplished by book cover assembly apparatus.\nIn a known cover fabrication apparatus, the cut-to-size covering material, or cover cloth, is separated from a magazine and fed to a glue roller via a cloth cylinder. A gripper bar acquires the glue-coated cloth and deposits it on a cover table. Feeder elements push the cut-to-size front and back cover boards from magazines to a ready-use supply station simultaneously with the infeeding of a spine board that has been cut to length from a reel.\nTwo cover boards and a spine board are simultaneously picked up at the supply station by a suction head mounted on an arm. The arm is rotated through 180.degree. to bring the boards into registration with a glue-coated cloth which is lying on the cover table. Due to the materials being brought into contact while substantially parallely oriented, i.e., the board and cloth are pressed flat against one another, attempts to increase production rate result in air being entrapped between the cloth and boards. Such entrapped air, in turn, results in blisters which give the book cover an unsightly appearance and, if they are comparatively large, will result in the cover being rejected by quality control."}
{"text": "1. Field of the Invention\nThe present invention relates to an intravascular catheter having a stiff proximal end portion and gradually softer and more flexible portions progressively toward a distal end of the catheter which is used for diagnostic, interventional and or drug infusion procedures such as blood clot dissolving drugs, chemotherapeutic agents, and injection of contrast media to visualize vasculature and or anatomy using a fluoroscope to assist visualization of the catheter inside the human body. The present intravascular catheter can also be used to deliver coils to aneurisms, and embolic agents to arteriovenous malformation (AVM) in the brain and other parts of the human vasculature such as uterine fibroids. The latter procedure is usually referred to as uterine fibroid embolization (U FE). The present invention more specifically relates to micro catheters for use in interventional neuro radiological procedures where access to brain vasculature is required. Further, the present invention relates to a microcatheter to deploy coils and embolic agents to treat certain brain vasculature syndromes to include but not limited to aneurisms and AVM indications.\n2. Description of the Related Art\nIntravascular catheters are used to diagnose and treat a number of medical conditions of the vascular system using a technique called angiography. A number of intravascular catheters are used to diagnose coronary artery disease related to stenosis and to determine hemodynamic factors such as cardiac output. In addition, smaller catheters, “micro-catheters”, are used to infuse certain blood clot dissolving drugs to the coronaries or to brain blood vessels in stroke related cases. Further, in addition, intravascular catheters are also used to deploy stents in the coronary arteries as well as in the peripheral vasculature; other indications may include but not limited to deployment of embolic agents and coils to target vasculature within the brain and elsewhere in the human body.\nIn order to properly navigate or manipulate a catheter after introduction in the human vasculature, it is imperative that the intravascular catheter be designed and constructed in such a manner as to facilitate introduction into a blood vessel and to support further manipulations to reach the target blood vessel. Manipulation of the catheter is done by the physician using the proximal segment of the catheter, after catheter introduction and external to the blood vessel; that is, the catheter needs to be rotated (torque) outside the body in the proximal segment of the catheter and a corresponding rotational reaction must be achieved in the distal tip segment inside the body. Therefore, a vascular catheter, to sustain torsional continuity from the proximal end to the distal end of the, while inside the vasculature, must transmit torque from the proximal end to the distal end inside the body. The typical distance between proximal and distal end for vascular catheters ranges between 80-120 centimeters and shorter for renal applications that require a length of approximately 45-55 centimeters; thus, a vascular catheter is required to advance forward when pushed, be able to be retrieved when pulled and to rotate when torque is applied as part of the maneuvering process to reach the target vessel. Another critical characteristic is that the catheter's wall must not collapse (kink) and retain its lumen integrity, cylindrical shape or ovality during the medical procedure.\nHeretofore, a number of analogous and non-analogous intravascular catheter constructions have been proposed, as disclosed in the following U.S. Patents:\nPat. No.Patentee3,485,234Stevens4,044,765Kline4,402,684Jessup4,430,083Ganz et al.4,516,972Samson4,581,390Flynn4,636,346Gold et al4,737,153Shimamura et al.4,842,590Tanabe et al.4,955,862Septka5,156,155King5,234,416Macauley et al5,279,596Castaneda et al.5,308,342Sepetka et al5,382,234Cornelius et al5,445,624Jimenez5,454,795Samson5,458,605Klemm5,554,139Okajima5,695,483Samson5,702,373Samson5,569,200Umeno et al5,782,809Umeno et al5,816,927Milo et al.5,879,342Kelly5,947,940Beisel5,951,539Nita et al.6,152,912Jansen et al.6,824,553Samson et al."}
{"text": "An electronic device includes a connector into which a cable for wired communication is inserted and a slot into which a removable storage medium is inserted in order to exchange information with another device. Therefore, the housing of an electronic device includes a hole for exposing an electronic device terminal such as a connector or slot to the outside of the housing. Hereinafter, in the present specification, a hole for exposing an electronic device terminal to the outside of the housing is referred to as a “socket”. Generally, the socket of an electronic device is covered and protected with a cover.\nRecent electronic devices are required to be waterproof and need to prevent water from entering through the sockets. Therefore, packings are provided on covers that cover the sockets to protect the electronic device terminals. Structures for preventing water from entering through the sockets are classified into a structure called a longitudinal compression type in which a packing is placed on the front side of the opening of the socket and a structure called a transverse compression type in which a packing is put in the socket and brought into contact with the wall. The longitudinal compression type needs to secure a packing margin around the opening of the socket, and thus hinders miniaturization of the electronic device, so the adoption of the transverse compression type is progressing.\nPatent Literature 1 discloses a transverse compression type waterproof structure."}
{"text": "Traditionally, contact information has been exchanged between parties through the use of business cards. Business cards often include contact information such as, for example, street addresses, telephone number(s), fax number, e-mail addresses and website addresses.\nTraditional business cards suffer from numerous disadvantages. In particular, traditional business cards may be lost and easily destroyed or damaged.\nAccordingly, there exists a need for improved methods and systems for exchanging contact information.\nSimilar reference numerals may have been used in different figures to denote similar components."}
{"text": "The invention is generally related to image processing systems and, more specifically, to a method and apparatus for performing geo-spatial registration within an image processing system.\nThe ability to locate scenes and/or objects visible in a video/image frame with respect to their corresponding locations and coordinates in a reference coordinate system is important in visually-guided navigation, surveillance and monitoring systems. Aerial video is rapidly emerging as a low cost, widely used source of imagery for mapping, surveillance and monitoring applications. The individual images from an aerial video can be aligned with one another and merged to form an image mosaic that can form a video map or provide the basis for estimating motion of objects within a scene. One technique for forming a mosaic from a plurality of images is disclosed in U.S. Pat. No. 5,649,032, issued Jul. 15, 1992, which is hereby incorporated herein by reference.\nTo form a xe2x80x9cvideo mapxe2x80x9d, a mosaic (or mosaics) of images may be used as a database of reference imagery and associated xe2x80x9cgeo-coordinatesxe2x80x9d (e.g., latitude/longitude within a reference coordinate system) are assigned to positions within the imagery. The geo-coordinates (or other image or scene attributes) can be used to recall a mosaic or portion of a mosaic from the database and display the recalled imagery to a user. Such a searchable image database, e.g., a video map, is disclosed in U.S. patent application Ser. No. 08/970,889, filed Nov. 14, 1997, and hereby incorporated herein by reference.\nA system that images a scene that has been previously stored in the reference database and recalls the reference information in response to the current images to provide a user with information concerning the scene would have applicability in many applications. For example, a camera on a moving platform imaging a previously imaged scene contained in a database may access the database using the coordinates of the platform. The system provides scene information to a user. However, a key technical problem of locating objects and scenes in a reference mosaic with respect to their geo-coordinates needs to be solved in order to ascertain the geo-location of objects seen from the camera platform\"\"s current location. Current systems for geo-location, the mapping of camera coordinates to the geo-coordinates, use position and attitude information for a moving camera platform within some fixed world coordinates to locate the video frames in the reference mosaic database. However, the accuracy achieved is only on the order of tens to hundreds of pixels. This inaccuracy is not acceptable for high resolution mapping.\nTherefore, there is a need in the art for a method and apparatus that identifies a location within an imaged scene with a sub-pixel accuracy directly from the imagery within the scene itself.\nThe disadvantages of the prior art are overcome by the present invention of a system and method for accurately mapping between camera coordinates and geo-coordinates, called geo-spatial registration. The present invention utilizes the imagery and terrain information contained in the geo-spatial database to precisely align the reference imagery with input imagery, such as dynamically generated video images or video mosaics, and thus achieve a high accuracy identification of locations within the scene. The geo-spatial reference database generally contains a substantial amount of reference imagery as well as scene annotation information and object identification information. When a sensor, such as a video camera, images a scene contained in the geo-spatial database, the system recalls a reference image pertaining to the imaged scene. This reference image is aligned very accurately with the sensor\"\"s images using a parametric transformation. Thereafter, other information (annotation, sound, and the like) that is associated with the reference image can easily be overlaid upon or otherwise associated with the sensor imagery. Applications of geo-spatial registration include text/graphical/audio annotations of objects of interest in the current video using the stored annotations in the reference database to augment and add meaning to the current video. These applications extend beyond the use of aerial videos into the challenging domain of video/image-based map and database indexing of arbitrary locales, like cities and urban areas."}
{"text": "This invention relates to liquid dispensing apparatus and methods, and more particularly to airless liquid dispensing nozzles and methods.\nLiquid dispensing systems include methods and apparatus using compressed air to atomize and shape a spray pattern for application to a substrate, and airless liquid dispensing systems in which liquid is forced through a nozzle, frequently at high fluid pressures, in an expanding fan-like sheet for atomization and application to a substrate, and also airless liquid dispensing systems in which liquid streams are directed from a nozzle for impingement and the formation of an expanding fan-like sheet in air.\nCompressed air spraying and dispensing systems, while providing flexibility in operation and variability in the shape and angle of a spray pattern through the adjustment of compressed air jets from a plurality of orifices in the dispensing nozzle, suffers a serious disadvantage because of liquid spray particles and vapors which are blown away from the substrate and into the operating environment, frequently in violation of regulations for safe operation and protection of the environment.\nAirless liquid dispensing systems suffer from a lack of flexibility during their operation. Airless liquid dispensing nozzles are designed to dispense liquid in a substantially constant and pre-selected angle of dispersion. Thus, if during operation it becomes desirable to change the width of liquid being dispensed from an airless operating system, it has been necessary to stop operation of the system, remove the nozzle being used, and replace it with a nozzle providing a more desirable angle of dispersion. This is inconvenient and time consuming, frequently requiring cleaning of the dispensing apparatus and nozzle.\nThe ability to vary the width of liquid being dispensed and applied to a substrate while continuing dispensing is of particular value in the application of plural component materials, such as polyesters in gel-coat, spray-up, and wet-out operations. In many such operations, the substrates and molds to which plural component materials, such as polyesters are being applied, present varied and complex shapes, frequently with channels and corners. The ability to apply, for example, polyester materials, in narrow widths to channels and corners, and in wide widths to larger, planar areas of a mold can assist the equipment operator in obtaining quick, complete and uniform coverage of the substrate or mold, and effective wet-out of reinforcing glass fibers or mat.\nThe invention provides an airless dispensing system in which liquids may be dispensed with variable included angles from substantially 0xc2x0 to as much as 50xc2x0-60xc2x0, by a simple adjustment made by a dispenser operator without ceasing operation or disassembling the dispenser apparatus.\nIn the invention, a variable angle liquid dispenser comprises a nozzle having a forward face with a dispensing orifice, a first passageway in the nozzle having a central axis intersecting the dispensing orifice, a plurality of angled second passageways in the nozzle, each angled second passageway having a central axis intersecting the dispensing orifice and the central axis of the first passageway, and a variable flow means adjustable to vary the flows of liquid entering the first passageway and the plurality of angled second passageways and to thereby vary the included angle of the liquid dispensed from the dispensing orifice.\nVariable angle liquid dispensers of the invention can comprise two major elements, a nozzle body forming the forward face, the dispensing orifice, the first passageway and the plurality of intersecting angled second passageways, and variable flow means comprising a body forming an input passageway and carrying a variable flow splitter between the input passageway and the passageways of the nozzle body. In one preferred embodiment, the variable flow means comprises an assembly including a flow divider body forming an input passageway leading to the first passageway and the plurality of angled second passageways of the nozzle, and a valve member movably carried by the body to provide a variable flow division between the first passageways and plurality of angled second passageways of the nozzle. A preferred flow divider body can have a forward portion that is adapted for sealed engagement with the nozzle body, can provide a seal between the first passageway and the plurality of angled second passageways, and can provide a first feed passageway between the input passageway and the first passageway of the nozzle body, and a second feed passageway with an entrance opening between the input passageway and the angled second passageways of the nozzle body. The valve member can be threadably and rotatably carried by the flow splitter body and can be moved variably with respect to the entrance opening of the second feed passageway to provide a variable flow division between the first feed passageway and the second feed passageway, thereby varying the portions of the liquid flowing in the input passageway that flow through the first passageway and angled second passageways of the nozzle.\nOne preferred nozzle body of the invention includes a plurality of passageways converging adjacent the dispensing orifice. A first passageway lies on the central axis of the nozzle body and two angled second passageways have their central axes lying outboard in the same plane as the central axis of the first passageway, and converging with an included angle of from about 25xc2x0 to about 50xc2x0 between them. The first passageway has a diameter about 25% to about 75%, preferably about 70% of the diameters of the angled second passageways and a central axis bisects the central axes of the two outboard passageways.\nIn one preferred method of the invention, flows of the two components of a plural component material, such as a polyester resin and a catalyst therefor, are directed, under pressure, from their sources for mixing and dispensing from a nozzle having a plurality of passageways, at least two of the nozzle passageways being angled, with their central axes converging at an included angle, for example, from about 25xc2x0 to about 50xc2x0, and with a central nozzle passageway bisecting the converging central axes of the two angled outboard passageways. The flow rates of the two components are controlled to provide desired flow rates for proper mixing, for example, a catalyst flow rate about 0.5% to about 10% of the flow rate of a polyester resin, and the two liquid components catalyst are mixed, while flowing, and the flowing mixed liquid components are directed to the plurality of passageways of the nozzle and are dispensed from the nozzle as combined and mixing streams, forming an expanding, substantially planar stream of further mixed two-component material for application to a substrate. The flow of mixed two-component material is variably divided between the at least two angled outboard passageways and the central passageway, and the included angle of liquid dispensed from the dispensing nozzle is varied from substantially 0xc2x0 to about 50xc2x0, varying the width of mixed plural component material applied to the substrate.\nOther features and advantages of the invention will be apparent from the attached drawings and more detailed description of the currently known best mode of the invention, which follows."}
{"text": "This invention relates to a process for the preparation of amino-substituted penicillins and, more particularly, to an efficient process for producing .alpha.-aminobenzylpenicillin and related ring-substituted compounds in high yield and purity.\n.alpha.-Aminobenzylpenicillin and .alpha.-amino-substituted-benzylpenicillins are well known in the art and numerous processes have been proposed for their production. In general these processes involve the reaction of 6-amino-penicillanic acid with an acylating agent such as the acid chloride, acid bromide, acid anhydride, and mixed anhydride of a derivative of .alpha.-aminophenylacetic acid or .alpha.-amino-substituted-phenylacetic acid in which the amino group is protected with a suitable protecting group.\nThese known methods for the preparation of .alpha.-aminobenzylpenicillins and .alpha.-amino-substituted benzylpenicillins by the acylation of 6-aminopenicillanic acid result in the preparation of mixtures which contain, in addition to the desired penicillin, unreacted starting materials, hydrolyzed acylating agent, and products of side reactions which are often difficult to separate from the desired penicillin reaction product.\nVarious ways have been proposed in the past to remove these unwanted materials such as, for example, the formation of insoluble arylsulfonic acid salts of .alpha.-aminobenzylpenicillin as described in U.S. Pat. No. 3,180,862, or the even more complex isolation process of U.S. Pat. No. 3,271,389, but in each case the known recovery procedures have been concentrated on removing the contaminants after the desired .alpha.-amino-substituted benzyl penicillins have been formed. Consequently, these known processes are characterized by complex procedures which at the very least increase the cost of producing the desired penicillin product."}
{"text": "1. Field of the Invention\nThis invention relates to seismic exploration in general. More specifically, it concerns an improved method of reflection-type seismic exploration that is particularly applicable to marine operations.\n2. Description of the Prior Art\nIt has been known for quite some time that in marine seismic operations, there are particular problems that are not encountered in land operations. One aspect stems from the fact that in marine operations the seismic detectors placed below the surface of the water are sensitive to seismic waves in the water regardless of their direction of travel. Furthermore, pressure-type detectors are usually used, whereas in land operations the detectors are ordinarily a displacements or inertia type.\nHeretofore, in marine operations, seismic waves were often generated by detonation of an explosive charge which was usually placed at a depth of ten feet or less below the surface. This avoided interference produced by the phenomenon commonly called \"bubble bounce\" which interference is generated by charges fired at greater depth. The shallow depth of charge also would avoid the problem of ghost reflections which are encountered in land-type shooting where the charge is detonated in competent earth material some distance below a good reflector.\nHowever, the additional problem remained in offshore seismic exploration which was created by the vertically travelling reflection signals that will reflect back down from the surface of the body of water and, consequently, will create an interference pattern.\nConsequently, it is an object of this invention to provide an improved method of marine-type seismic surveying. Furthermore, another object is to provide a method of marine-type seismic exploration which greatly simplifies the procedure for elimination of undesirable seismic wave energies from a recording.\nAnother object of the invention is to provide a marine reflection-type seismic method for creating a record which is independent of depth of the detectors in the water and free of down-going signals."}
{"text": "1. FIELD OF THE INVENTION\nThe present invention relates to a semiconductor device and a method of manufacturing the same, and more particularly, a thin film semiconductor device formed on an insulating substrate and a method of manufacturing the same.\n2. DESCRIPTION OF THE RELATED ART\nSemiconductor devices such as transistors can be formed on an insulating substrate such as a glass plate or a silicon substrate coated with a passivation layer. It is not possible to apply epitaxial growth onto the insulating substrate in a usual manner. Thus, a thin film semiconductor layer of polycrystalline or amorphous state is first formed on the insulating substrate and a semiconductor device is then formed on the thin film semiconductor layer.\nFor example, a matrix of cells together with active elements such as transistors are formed on a transparent substrate, thereby forming an active matrix. In the field of liquid crystal displays, a high attention is directed to an active matrix liquid crystal display panel in which a matrix of transparent electrodes for exciting a liquid crystal material and thin film transistors (TFTs) serving as switching elements are formed on a glass substrate. Such a TFT can be made in the form of either a p-channel insulated gate (IG) or an n-channel IGTFT. In the case of making the n-channel IGTFT, n.sup.+ -type source/drain regions are formed by doping a thin film of highly resistive i-type or p-type silicon (of polycrystalline or amorphous state) with an n-type impurity, and source/drain electrodes are disposed on the n.sup.+ -type source/drain regions in contact therewith. In the case of making the p-channel IGTFT, p.sup.+ -type source/drain regions are formed in a thin film of i-type or n-type silicon with a p-type impurity, and source/ drain electrodes are disposed on the p.sup.+ -type source/drain regions in contact therewith.\nLepselter and Sze have proposed an IGTFT using Schottky barrier contacts for a source and a drain (see Proceedings of the IEEE, Proceedings Letters, Aug. 1968, pp. 1400-1402). In the proposed IGFET, source and drain electrodes made of platinum silicide PtSi are brought in contact with a &lt;100&gt; oriented n-type silicon substrate having a resistivity of 1 .OMEGA..multidot.cm. When a negative voltage is applied to a gate electrode, a channel inverted to be of p-type is produced between the source electrode and the drain electrode. In the case where the electrodes of PtSi are disposed on the n-type silicon substrate, the barrier height of 0.85 V is established. But, in the case of electrodes of PtSi disposed on a p-type silicon substrate, the barrier height is 0.25 V so that a drain current flows.\nTaking account of the fact that when the source and drain regions are formed with Schottky junctions, a short-channel effect can be improved but a restriction is imposed on an ability of current supply from the source region, Mizutani has proposed an MOSFET having a source region formed through impurity diffusion and a drain region formed with a Schottky junction (see JP-A-58-182871).\nFor driving the liquid crystal display (LCD) panel, there are required various peripheral circuits including a shift register, a matrix circuit, an inverter circuit, etc. If it is possible to incorporate these peripheral circuits into the LCD panel, the number of parts required can be reduced, thereby allowing improved reliability and greatly lowered cost.\nWhen it is desired to make a limitation of power consumption , complementary insulated gate thin film transistors (C-IGTFTs) must be used for a part of such peripheral circuits. Namely, it is necessary to simultaneously fabricate a p-channel IGTFT and an n-channel IGTFT. The fabrication of C-IGTFTs requires both doping with an n-type impurity and doping with a p-type impurity. This increases the number of process steps for manufacture of the device, which provides a great factor of increasing the cost. For example, in spite of the fact that switching elements necessary for ah active matrix type of LCD panel can be formed by only either p-channel IGTFTs or n-channel IGTFTs, the requirements for incorporation of the peripheral circuits into the panel necessitate the use of both p-channel IGTFTs and n-channel IGTFTs and hence the doping with both p- and n-type impurities, thereby resulting in remarkable increase of the number of photomasks to be used and of the number of process steps to be carried out."}
{"text": "Electronic content, such as web pages, search results, and other types of documents, often includes links to other documents, web pages, and the like. A link is a connection from one page to another that can be selected from a first page to cause the other page to appear in a web browser application or the like. The links on a page can be defined by the author of the page, or can be added to the page automatically, e.g., by an advertising system that adds advertisements to the page. Pages can be generated by an automated system such as a search engine, which identifies web pages that are available via a network such as the Internet, and adds links to those pages based upon the addresses that are found by the search engine. Attracting users to a web site can be desirable for a number of reasons. Illegitimate or malicious web sites can to attempt to gather private information such as email addresses and passwords from users, e.g., through phishing attacks. There are also more benign reasons to attract users to a web site, such as to increase the site's traffic, the number of times advertisements have been viewed on the site, and so on.\nWeb link spoofing attacks have been developed to attract users to web sites that the users do not intend to visit. These spoofing attacks deceptively present an illegitimate web link that appears legitimate. For example, suppose that a web site named Good Web Site has a legitimate web link good.com. The link good.com looks similar to the link g00d.com, in which each letter o is replaced by the number 0. The two links look particularly similar if they are displayed in uppercase, i.e., G00D.COM and GOOD.COM. As another example, the letter l in a legitimate link can be changed to a number 1 to create an illegitimate link that is visually similar to the legitimate link.\nURL's can contain characters from numerous international languages. There are a number of characters in different languages that look alike. Characters that look alike are referred to as homographs. URL spoofing attacks that take advantage of the visual similarities between different characters that can be from different languages are thus referred to as Internationalized Domain Name (IDN) homograph attacks. For example, the English letter c (pronounced cee) looks similar to the Russian letter c (pronounced ess). A URL that includes an English c, such as chase.com, can be spoofed by a URL that uses a Russian c in place of the English c, and looks very similar, such as chase.com. A user can be lured to an illegitimate version of the chase.com web site that is registered to the spoofed chase.com domain name by presenting a hyperlink having a URL that refers to the spoofed chase.com. Users are unlikely to see the difference between the legitimate and spoofed domain names, and thus unlikely to be aware that they are accessing an illegitimate web site, particularly if the illegitimate web site's appearance is similar to that of the legitimate chase.com site.\nBecause of such visual similarity between different characters, users can be lured into clicking on or selecting the illegitimate link when they intend to access the legitimate web site. When the user follows the web link to the illegitimate g00d.com web site, the illegitimate web site is loaded and displayed, and the spoofing attack has succeeded. The illegitimate site can, for example, display information or advertisements, attempt to convince the user to perform a transaction, request information from the user, attempt to install malware or spyware on the user's computer, and perform other malicious or potentially damaging operations. Spoofed web links can lead users to phishing attacks, in which an illegitimate site is designed to mimic sites that contain important user information and convince the user to login, thereby providing an attacker with their user name and password.\nA web link ordinarily has two parts: link text and a reference to a target web page, such as a Uniform Resource Locator (URL) that identifies the target web page. The link text is displayed on a web page to visually represent the link. The link text can be clicked on or selected to cause a web browser to load the target web page referred to by the URL. The link text can also be referred to as anchor text, a link label, or a link title.\nFor example, in the Good Web Site example, a link to the site can have link text such as “Good Web Site” or “www.good.com”, and a link URL such as www.good.com. In one aspect, a legitimate link URL correctly references the web site described or implied by the link's link text, such as www.good.com. An illegitimate link has a URL, such as g00d.com, that references a web site different from that described or implied by the link text. The link text is not necessarily the same as the link URL, and can be a description or name of the web page referred to by the link instead of a textual copy of the link URL. However, illegitimate links often set the link text to a URL, at least in part because users are more likely to trust and follow a link that is displayed as a legitimate-looking URL, as opposed to a link displayed as a word or phrase. Therefore, illegitimate links can set the link text to a legitimate URL and set the URL to an illegitimate link in an attempt to lure users into following the legitimate URL. Alternatively, illegitimate links can set the link text to an illegitimate URL, e.g., g00d.com, that looks similar to the legitimate URL, such as good.com, and again set the link URL to the illegitimate URL, g00d.com, so that a comparison of the characters that represent the link text to the characters that represent the link URL will indicate that both are the same, and such a comparison will not identify the link as illegitimate.\nAlthough the link text is ordinarily displayed on web pages to represent the link, the user is able to view the link URL itself, e.g., by placing a cursor or mouse pointer over the link text, and when a user actually opens the illegitimate page. Thus, the user can then attempt to visually verify that the URL is legitimate by placing the mouse pointer over the link text prior to clicking on the link, and checking the URL that is displayed. The user can also attempt to visually verify the URL by clicking on the link, allowing the target page to begin loading, and visually verify the URL that is displayed in the browser's address bar. In either situation, if the URL appears to be illegitimate, e.g., because it references a web site or contains text that does not appear to be related to the link text, then the user can decide to ignore the link or the loaded target page. However, if the URL appears to be legitimate, then the user is likely to follow the illegitimate link or read the illegitimate loaded target page. It would be desirable, therefore, to protect users against web link spoofing, so that users do not unintentionally access illegitimate web pages.\nExisting techniques for blocking web link spoofing attacks include filtering based on heuristics, and blocking sites that appear on lists of known unsafe pages. The heuristics can be used to identify suspicious messages and, and require additional effort, e.g., a confirmation input, by users. Both the filtering and the site blocking lists can fail against modern attacks. The filtering technique can fail for a number of reasons, such as a relatively high false-positive rate that leads users to disable the features or ignore warnings, even for content that is actually an attack. Further, the attacker can design messages to avoid detection, e.g., by indirectly determining the heuristics used by the filter, or by directly testing their messages against their own copy of the filtering software. For example, an attacker could send their message to themselves, and change the message until it passes through the filtering software. The site blocking lists fail because of the delay between the start of the attack and detection of the attack. A potentially large number of victims can be attacked before the attack is detected. Neither of these techniques works against attacks that are specifically targeted against a small number of victims. Targeted attacks involve tailoring messages to bypass filtering, and the small volume of attacks reduces the likelihood of the phishing site being detected at all, let alone early enough to detect all attacks."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for positioning adapted for use in an exposure apparatus for manufacturing semiconductor devices, and more particularly to a method and an apparatus, in positioning a substrate provided with at least a registration mark and other patterns, for positioning by distinguishing said mark from said other patterns.\n2. Description of the Prior Art\nWith the progress toward finer patterns in semiconductor devices, particularly in large-scale integrated circuits, there are commonly employed reduced projection exposure apparatus for circuit pattern printing in order to meet the requirements of fine pattern definition and a high productivity. Such conventional apparatus projects a reticle pattern, which is larger for example 5 times than the pattern size desired on the silicon wafer, in a reduced size through a projection lens, thus exposing a square area with a diagonal length of 21 mm or shorter on the wafer in one exposure. Consequently, in order to print the circuit patterns on the entire surface of a wafer of a diameter of ca. 125 mm, there is employed a so-called step-and-repeat process in which the wafer is placed on a movable stage and is subjected to repeated exposures with stepwise movement.\nIn the manufacture of large-scale integrated circuits, there are formed patterns of at least several layers on a wafer in succession, and the desired function cannot be attained due to deficient conductivity or insulation unless the error in superposition or positional aberration between the patterns of different layers is maintained under a determined limit. For example, in a circuit with a minimum line width of 1 .mu.m, there is only permitted a positional aberration of 0.2 .mu.m at maximum.\nIn the reduced projection exposure process, the pattern registration, namely the registration of a projected reticle pattern with a pattern already formed on the wafer, is achieved either by the off-axis method or by the through-the-lens method. In any case, the wafer is subjected to a rough alignment, called pre-alignment, when placed on the stage, and, if said pre-alignment is not precise enough, the registration mark formed on the wafer is significantly displaced from the central detecting position in an alignment microscope in case of the off-axis method or in a detecting system in case of the through-the-lens method. This leads to an erroneous detection of another pattern as the registration mark and eventually gives rise to an error in registration."}
{"text": "By this invention, I have reduced the size, complexity, and expense of equipment for treating liquid with ozone on a demand basis; and by carefully selecting and combining components, I have been able to make smaller scale ozone treatment equipment operate conveniently and reliably for safely treating liquid rapidly for small volumes and flow rates. Equipment according to my invention can be operated under a residential countertop, for example, to treat liquid flow on a small scale."}
{"text": "For the ignition of briquets and charcoal for grills up till now ignescent fluid has been the dominating and sole accepted ignition aid for producing in an acceptably short time embers for broiling. Among the drawbacks of ignescent fluid are the hazards of the ignition procedure. At times, ignescent fluid has been confused with other fluids and caused severe burns in children and in some known instances children have been poisoned by drinking the fluid. In addition, the ignescent fluid is bulky and generally difficult to bring along. It sometimes also imparts obtrusive flavours to the food being broiled. The use of ignescent fluid is also expensive.\nTo light a fire in fireplaces, furnaces, and suchlike, one normally uses newspaper leaves and the like, in conjunction with wood chips. This is a time-consuming method. Ignition aids known as `fire lighters` may also be used. A method for producing fire lighting aids was described in SE-A-No. 41 897, in 1914. According to this method, paper, sulphite or sulphate pulp, is impregnated with a combustible substance which is either liquid or solid, such as resin, resin dissolved in some combustible substance such as spirits, turpentine, raw or refined petroleum, tar, or some other suitable substance. After being impregnated, the paper or the pulp is rolled onto spindles, and fire lighting aids then prepared from the strips, whether wet or dry, the final product being in the form of small reels. According to SE-A-No. 96 174 fire lighters are produced from lumbering or wood mill debris, which is cut into chips, defibrated, mixed with water to achieve a suitable consistency and lastly formed into a plate, which is dewatered by pressing and then dried. This plate is dipped in molten paraffin, stearin, or tallow or a mixture of these at a temperature of 80.degree.-100.degree. C. After drying, the plate is cut into pieces of a certain width and length. Before being impregnated, the plate is provided with grooves, to facilitate the cutting of the plate into small square blocks.\nA drawback which is common to these and other known fire lighting aids is that the area of combustion is small, the product thus having to be ignited at a very small area. Therefore, it is not at all uncommon to fail at the ignition of these products, even if the burning time may be long. In addition, the positioning of the lighter is critical, for instance when lighting a fire on a grill, since the lighter, being very small, may easily fall down through the grid.\nAnother known lighting aid consists of cubes of a brittle material which easily crumbles and has a strong odour, so that the product must be carefully packed and gently handled.\nParaffin impregnated cellulose pulp is a better lighting aid. The area of combustion of this product in relation to its volume is greater, and hence the product burns more intensely and over a larger area. Even though its burning time is shorter than that of a more compact product of the same volume, the fire or the bed of briquettes or coal to be ignited is lit more effectively and more safely. Another desirable property of the lighter is that it is free of tackiness. Nor should it crumble when broken, as is the case if not all paraffin has become absorbed into the pulp. At the same time it must be water-repellent and inflammable. These demands have caused considerable manufacturing problems."}
{"text": "A process of color separating a transmitted light or reflected light from a color original into so-called three primary colors of R (red), G (green), and B (blue) by an ordinary optical method using, for example, red, green and blue-purple color filters, making negative or positive films for blue printing (cyan), red printing (magenta), and yellow printing (yellow) basing on the lights of three colors of R, G, and B, and then making each color separated printing plate using each film thus obtained has hitherto been widely performed. In the process, photographic light-sensitive materials having light-sensitive characteristics in the whole wavelength region of visible light of from 400 n.m. to 700 n.m. are required.\nHitherto, as the light-sensitive materials having such light-sensitive characteristics, silver salt light-sensitive films are used. However, the silver salt light-sensitive films are expensive as compared to ordinary films and in the case of using such silver salt light-sensitive films, it is difficult to handle the films in bright room. In the case of using a silver salt light-sensitive material, it cannot be directly used as a printing plate but a process of developing and fixing the silver salt light-sensitive film and then printing the positive or negative images formed on the silver salt film onto an original printing plate such as a presensitized printing plate (PS printing plate), etc., to make a printing plate is required.\nThese sequential operations per se, are troublesome and require specific apparatuses for each operation. An expense for these apparatuses, and personnel, as well as material cost for light-sensitive films cannot be disregarded.\nThis invention has been achieved in view of forgoing problems to provide color separating printing plate directly from the color originals in a short time."}
{"text": "1. Field of the Invention\nThe present invention generally relates to pumps, and more particularly to positive-displacement rotary pumps.\n2. Description of the Related Art\nPositive displacement pumps displace a known quantity of liquid with each revolution of the pumping elements (e.g., vanes). Positive displacement pumps displace liquid or gas by creating a space between the pumping elements and trapping the liquid or gas within the space. Rotation of the pumping elements then reduces the volume of the space and moves the liquid out of the pump. A rotary vane pump is an example of a positive-displacement pump.\nRotary vane pumps operate through the action of a number of rotating vanes or blades. A conventional rotary vane pump includes a rotor assembly eccentrically positioned within a pumping chamber. The number of vanes are spaced around the rotor to divide the pumping chamber into a series of cavities. As the rotor rotates, these cavities rotate around the pumping chamber continually changing in volume due to movement of the vanes and the eccentric alignment of the rotor and pumping chamber. An inlet communicates with the pumping chamber on the side of the pump where the volume of the cavities expand. Similarly, an outlet communicates with the pumping chamber on the side of the pump where the volume of the cavities contract. As each cavity expands, a partial vacuum is created to draw fluid into the pump through the inlet. As the cavity contracts, the pressure within the cavity increases forcing the fluid out of the pump through the outlet. This expansion and contraction process continues for each cavity to provide a continuous pumping action.\nThere is a desire to improve upon the currently available rotary pumps. For example, there is a desire to reduce the cost of manufacturing rotary pumps while maintaining (and possible increasing) the vacuum level produced by a pump of specific dimensions. There is also the desire to increase the volume of fluid that can be displaced during a period of time by a pump of specific dimensions (i.e., without increasing the overall dimensions of the pump). Further, there is the desire to simplify the manufacturing and assembly required for producing rotary pumps.\nThe present invention is directed to a dual chamber or double sided rotary pump that includes a stator housing and a rotor.\nIn accordance with an embodiment, the stator housing has an oblong inner surface. The rotor, which is disposed in the stator housing, has a substantially circular outer surface within which a plurality of vane slots are defined. A first chamber is defined between a first half of the oblong inner surface and the outer surface of the rotor. Similarly, a second chamber is defined between a second half of the oblong inner surface, diametrically opposite the first half, and the outer surface of the rotor. Resting within each of the plurality of vane slots is a corresponding sliding vane. A first inlet port and a first outlet port provide access to the first chamber. Similarly, a second inlet port and a second outlet port provide access to the second chamber. The vane slots are arranged about the outer surface of the rotor such that there is always at least one of the vanes separating each of the first inlet port, the first outlet port, the second inlet port and the second outlet port from one another.\nAs the rotor is rotated within the stator housing, centrifugal force pushes or urges the vanes radially outward against the inner surface of the stator housing. As this occurs, each of the first and second inlet ports draws in fluid (i.e., gas and/or liquid), and each of the first and second outlet ports expels fluid. More specifically, fluid drawn into the first inlet port is expelled out of the first outlet. Similarly, fluid drawn into the second inlet port is expelled out of the second outlet port. This occurs as described below.\nAt any given time there exists multiple cavities formed between adjacent pairs of the vanes. For example, there are eight cavities in the embodiment of the present invention where there are eight vane slots and eight vanes. During each full rotation of the rotor, each formed cavity expands and contracts in volume twice. More specifically, each cavity expands in volume as it passes the first inlet port, shrinks in volume as it passes the first outlet port, expands in volume as it passes the second inlet port, and shrinks in volume as it passes the second outlet port. When a cavity expands in volume it creates a partial vacuum, as it passes one of the inlets ports, and thereby draws fluid into the cavity. When the same fluid filled cavity shrinks in volume, as it passed one of the outlet ports, it expels that fluid. Thus, at any given time (while the rotor is rotating at a sufficient speed) two chambers are drawing fluid in and two other chambers are expelling fluid. The remaining chambers are in the process of transferring fluid that has just be drawn in (by one of the input ports) toward one of the outlet ports, so that the fluid can be expelled.\nThe rotary pump further includes first and second side plates (also referred to as end caps) located opposite one another at axial ends of the stator housing. The first and second side plates together with the stator housing form a hollow oblong cylinder within which the rotor is disposed. One of the side plates may be integrally formed with the stator housing.\nIn accordance with an embodiment of the present invention, most or all of the rotary pump is manufactured out of plastic. This can significantly reduce the cost and weight of the rotary pump. In accordance with an embodiment, the stator housing and side plates are manufactured from polyetherimide, the rotor is manufactured from polyphenylene sulfide, and the vanes are manufactured from thermoplastic polyimide. For strength, durability and lubrication: the polyethermide can include a carbon fill of about 25-35 percent and a polytetrafluoro ethylene fill of about 10 to 20 percent; the polyphenylene sulfide can include a carbon fill of about 35-45 percent; and the polyimide can include a carbon fill of about 25-35 percent and a polytetrafluoro ethylene fill of about 10 to 20 percent.\nFurther embodiments, features and advantages of the present invention may be more readily understood by reference to the following description taken in conjunction with the accompanying drawings and claims."}
{"text": "1. Field of the Invention\nThis invention relates to hard floor panels.\n2. Related Technology\nIn the first instance, the invention is intended for so-called laminated floors, but generally it can also be applied for other kinds of floor covering, consisting of hard floor panels, such as veneer parquet, prefabricated parquet, or other floor panels which can be compared to laminated flooring.\nIt is known that such floor panels can be applied in various ways.\nAccording to a first possibility, the floor panels are attached at the underlying floor, either by glueing or by nailing them on. This technique has a disadvantage that is rather complicated and that subsequent changes can only be made by breaking out the floor panels.\nAccording to a second possibility, the floor panels are installed loosely onto the subflooring, whereby the floor panels mutually match into each other by means of a tongue and groove coupling, whereby mostly they are glued together in the tongue and groove, too. The floor obtained in this manner, also called a floating parquet flooring, has as an advantage that it is easy to install and that the complete floor surface can move which often is convenient in order to receive possible expansion and shrinkage phenomena.\nA disadvantage with a floor covering of the above-mentioned type, above all, if the floor panels are installed loosely onto the subflooring, consists in that during the expansion of the floor and its subsequent shrinkage, the floor panels themselves can drift apart, as a result of which undesired gaps can be formed, for example, if the glue connection breaks.\nIn order to remedy this disadvantage, techniques have already been through of whereby connection elements made of metal are provided between the single floor panels in order to keep them together. Such connection elements, however, are rather expensive to make and, furthermore, their provision or the installation thereof is a time-consuming occupation.\nExamples of embodiments which apply such metal connection elements are described, among others, in the documents WO 94/26999 and WO 93/13280.\nFurthermore, couplings are known which allow coupling parts to snap fit into each other, e.g., from the documents WO 94/1628, WO 96/27719 and WO 96/27721. The snapping-together effect obtained with these forms of embodiment, however, does not guarantee a 100-percent optimum counteraction against the development of gaps between the floor panels, more particularly, because in fact well-defined plays have to be provided in order to be sure that the snapping-together is possible.\nFrom GB 424.057, a coupling for parquetry parts is known which, in consideration of the nature of the coupling, only is appropriate for massive wooden parquetry.\nFurthermore, there are also couplings for panels known from the documents GB 2.117.813, GB 2,256.023 and DE 3.544.845. These couplings, however, are not appropriate for connecting floor panels."}
{"text": "1. Technical Field\nThe present invention relates to a fuel cell and to a method of manufacturing the fuel cell.\n2. Description of the Related Art\nIn a fuel cell, the output voltage that can be provided by an individual cell may be a fixed value, determined by the electrochemical reactions. As such, in order to supply the operating power required by an electronic product, it may be necessary to use a DC-DC converter or serially connect the individual cells to increase the output voltage.\nAccording to the method of arranging the individual cells for connection, a fuel cell can be divided into a bipolar stack or a monopolar stack.\nFIG. 1 is a cross-sectional view of a fuel cell according to the related art, in which a bipolar stack structure is shown. In the case of medium- or large-sized fuel cells, most of the fuel cells take the form of a bipolar stack. As illustrated in the drawing, cells which include anodes 2, 3 and cathodes 2, 4 formed in contact with a membrane 1, and separation plates may be stacked alternately in layers.\nThe bipolar stack structure may entail a large volume for the fuel cell, and are thus applied more often in medium- to large-sized fuel cells. Since the power supply used in portable electronic equipment may desirably be given a thin shape, it can be problematic to employ a fuel cell of a bipolar stack structure for such a power supply, due to the large thickness involved. As such, there can be problems in implementing a compact size and in providing a high output.\nTo overcome these problems of the bipolar stack structure, a fuel cell can be implemented to have a monopolar stack structure. The fuel cell of a monopolar stack structure can provide a higher output density per volume, and enables the supply of fuel without external power, so that the form of the fuel cell may be varied relatively freely. As such, the monopolar stack structure is often employed in small-sized fuel cells.\nThe monopolar stack structure can in turn be divided into a banded structure or a flip-flop structure.\nThe banded structure may require a connection crossing the membrane, while the flip-flop structure may be of a simple shape, having only a single-layer connection.\nWhile the banded structure may simplify the method of supplying fuel and air, the banded structure may require additional complicated equipment for connecting the electrodes in serial.\nWith the flip-flop structure, it is possible to naturally form a serially connected stack by having adjacent unit cells share a common electrode. However, the fuel cell thus formed may entail an extremely complicated flow path for supplying the fuel and air."}
{"text": "1. Field of the Invention\nThis invention relates to an extraction cleaning machine and, more particularly, to an upright extraction cleaning machine.\n2. Description of Related Art\nUpright extraction cleaning machines have been used for removing dirt from surfaces such as carpeting, upholstery, drapes and the like. The known extraction cleaning machines can be in the form of a canister-type unit as disclosed in U.S. Pat. No. 5,237,720 to Blase et al. or an upright unit as disclosed in U.S. Pat. No. 5,500,977 to McAllise et al. and U.S. Pat. No. 4,559,665 to Fitzwater.\nCurrent upright extraction cleaning machines can be made easier to use by limiting the weight and number of components, such as fluid storage tanks, on the pivoting handle of the upright cleaning machine. Reducing the weight that a user must support as the handle is tilted rearwardly can also lower the center of gravity for the machine, which results in a better feel to the user.\nFurthermore, the current extraction cleaning machines can be made easier to use and better adapted for a variety of cleaning conditions. For example, none of the current extraction cleaning machines includes an elevator responsive-to-handle position for restraining a floating roller-type agitation brush, which is automatically height adjustable in response to changes on the surface being cleaned. Another problem inherent with the known extraction cleaning machines is the difficulty of filling and emptying the fluid supply chamber and fluid recovery chamber, particularly with bladder-within-a-tank type assemblies. Further, none of the current upright extraction cleaning machines are simply convertible to a pre-spray applicator for directing cleaning solution to and agitating the surface to be cleaned without applying suction. Finally, current extraction cleaning machines do not use a the same motor to drive an agitation brush as well as an impeller. Is some cases a separate motor is used. In other cases, a turbine is used to drive the agitation brush or brushes which diminishes the suction power available to extract the dirty solution from the floor surface.\nA more recent development in the extraction cleaning industry is the use of steam or hot water as a cleaning agent. The cleaning machine incorporates a boiler or other means for generating steam or hot water, which is pumped to an applicator where it is brought into contact with the surface being cleaned. Because the steam is airborne, it may be unsafe to include detergents and the like in the cleaning solution. Further, while the steam systems have the advantage of creating a temperature that effectively kills a wider range of microbes, bacteria, microorganisms, and mites, the steam systems generally suffer from poor cleaning performance. Additionally, the high power requirement for generating steam may not be sufficient with ordinary 120V power supplies for running a vacuum motor as well as the steam generator, so cleaning performance is further hindered. Also, by adding a heater to a fluid supply chamber, the user may be inconvenienced by the amount of time required to heat the contents of the supply chamber to the desired temperature. Conversely, conventional detergent cleaning systems are somewhat effective at cleaning surfaces, but could be made more effective by raising the temperature of the cleaning solution to some temperature below the boiling point. There is an optimal temperature at which cleaning performance is maximized without causing damage to carpets or setting stains. This temperature is around 150.degree. Fahrenheit."}
{"text": "1. Field of the Invention\nThe present invention relates to nanoimprint technology, and more particularly, to a method of fabricating a nanoimprint mold which is essential to nanoimprint technology.\n2. Description of the Related Art\nRepresentative nanolithographic techniques for creating nanostructures with dimensions of less than 100 nm include electron beam (E-beam) lithography (EBL), focused ion beam lithography (FIBL), nanoimprint lithography (NIL), and deep UV (DUV) lithography. Recently, next-generation lithographic techniques such as extreme UV (EUV) lithography, X-ray lithography, hologram lithography, and laser interference lithography (LIL) have been developed.\nEBL or FIBL with a resolution of several nanometers (nm) has disadvantages in that accelerated electrons or ions are greatly affected by the conductivity or structure of a substrate, and throughput is low due to raster or vector scanning.\nDUV or EUV lithography which is mainly used in a semiconductor fabrication process and can generate wafer-scale patterns over a large area in a single step has disadvantages in that manufacturing equipment is expensive, costs of a photomask are high, and it is difficult to use a flexible substrate or material which is affected by a developer used for pattern formation after exposure.\nX-ray lithography which is considered as next generation lithography can provide a diffraction limited resolution of several nanometers, which is higher than that of optical lithography and similar to that of the EBL or FIBL. However, the X-ray lithography is not suitable for forming patterns over a large area because it uses a very large synchrotron as an X-ray source and it is difficult to increase the size of the X-ray source.\nLIL which is often used in forming repetitive patterns such as gratings can be easily applied to a large area, is rarely affected by the characteristics of a substrate, and can be achieved at low cost. Despite the advantages, the LIL cannot be widely used because it is difficult to form various patterns and requires strict overlay accuracy while performing a process of forming multi-layers.\nImprint lithography for forming patterns using a mold is largely divided into a thermal method and a UV curing method. The imprint lithography can be applied to a general semiconductor substrate, such as a silicon substrate, a plastic substrate requiring a low temperature process, and also to a non-conductive glass or quartz substrate.\nRecently developed nanoimprint lithography has become an attractive alternative to hot embossing lithography because of its easy control of pressure or temperature, and particularly, UV nanoimprint lithography can produce patterns even at room temperature and under normal pressure.\nHowever, the UV nanoimprint lithography requires a transparent substrate or mold which is generally formed of quartz. It is difficult to form sub-100 nm fine patterns over a large area of the quartz mold or substrate, thereby increasing manufacturing costs."}
{"text": "1. Field of Invention\nThe present invention is directed to a lubricant system for medical devices. More particularly, the present invention is directed to a non-volatile lubricious coating for hypodermic needles, catheters, and the like.\n2. Description of Related Art\nIt has become commonplace in the medical field to provide medical devices with lubricants for ease of use. For example, hypodermic needles are widely used in delivering and withdrawing fluids in medical practice. As originally used, hypodermic needles were used many times, the needles being sterilized between usages. A practitioner would sharpen the needles when they became dull, and then sterilize them prior to the next usage. Since the needles were reused, and often may have needed sharpening, the presence or absence of any lubrication on the outer surface of the needle had little effect on the penetration force or the pain perceived by the patient who was the recipient of the needle. With the development of commercially manufactured disposable needles that always have a fresh well-sharpened point, there was recognition that lubrication of the needle substantially reduced the pain perceived by the patient when a needle was administered to them.\nA convention is followed in this disclosure wherein the portion of a device toward the practitioner is termed proximal and the portion of the device toward the patient is termed distal.\nA tissue penetration by a hypodermic needle involves a sequence of events that collectively are perceived by the patient as whether or not the penetration causes pain. A distal point of the needle first touches the skin surface, stretches it, the point then cuts into the surface and begins penetration into the tissue. As the shaft of the needle passes through the original cut and into the tissue, there is also sliding friction of the tissue against the needle surface. In the hypodermic needle art when the forces for performing a hypodermic needle penetration are measured, the force measured prior to the needle point cutting the tissue is termed the “peak penetration force”, also called “F2” and the force required to continue the penetration into the tissue is called the “drag force” or “F4”. One primary component of the drag force is the sliding friction of the tissue against the surface of the needle shaft.\nInsertion of intravenous (IV) catheters into a patient causes similar issues regarding ease of insertion and patient discomfort. For example, IV catheters are designed to infuse normal intravenous solutions, including antibiotics and other drugs, into a patient. These catheters are also used to withdraw blood from the patient for normal blood-gas analysis as well as other blood work. The most common type of IV catheter is an “over the needle” catheter, in which a catheter is disposed over an introducer needle or cannula, which is used to insert the IV catheter into a patient. The needle is typically stainless steel and is hollow. Its distal tip is ground to a sharp tip for easy insertion into the patient. The catheter is also hollow and is disposed such that the sharpened tip of the needle is extended from the catheter for piercing of the patient's skin during use. Once the skin and vein have been pierced, the catheter is advanced over the needle and the needle is removed from the catheter. The catheter is typically extruded out of suitable plastic material such as TEFLON material (polytetrafluoroethylene), polyvinyl chloride, polyethylene, polyurethane or polyether urethane.\nThe use of lubricants on the surface of such hypodermic needles and IV catheters significantly reduces both the peak penetration force and the drag force. As a result, almost all single-use sterile disposable needles and IV catheters are supplied with a lubricant already applied to substantially the entire outside surface. A number of lubricants have been developed for use in such applications. Typically, such lubricants involve a medical grade polydimethylsiloxane which is commonly applied to the surface through a volatile carrier solvent which rapidly evaporates. For example, U.S. Pat. No. 5,911,711 to Pelkey discloses a lubricant system for hypodermic needles which includes a first layer formed from an at least partially cured organosiloxane copolymer and a polydimethylsiloxane that has a viscosity greater than about 1000 centistokes, and a second layer over the first layer that includes a polydimethylsiloxane having a viscosity of 50-350 centistokes. The coating compositions of the first and second layers are applied through a volatile carrier solvent such as a chlorofluorocarbon (CFC), and the first layer is thermally cured by applying heat. Unfortunately, volatile solvents such as CFC's raise significant environmental concerns.\nThere is a need in the medical industry for lubricants for medical devices such as catheters and needles which are environmentally friendly, which are easy to apply, and which do not involve the use of volatile organic solvents such as CFC's."}
{"text": "1. Field of the Invention\nThe invention relates to a structure for welding airport concrete steel mesh, wherein each separate transversely extending reinforcing rod in the steel mesh is simultaneously welded to all the longitudinally extending spaced apart reinforcing rods, the separate transversely extending rods are automatically fed onto the longitudinally extending reinforcing rods of a welding position, and the welded transverse and longitudinal rods are automatically indexed away from the welding position.\n2. Description of the Prior Art\nPrior automatic mesh welders have usually not been capable of welding heavy reinforcing rods such as 3/4\" reinforcing rods into steel mesh such as that required for airport concrete runway reinforcement. Further, wherein steel mesh welding has been accomplished in the past, welding has been limited as to spacing between the reinforcing members due to welding gun diameter limitations, and welds have not always been satisfactory due to low welding gun pressures available. Wherein satisfactory mesh welding has been accomplished in the past, the equipment necessary therefor and the methods used have been complicated and therefore uneconomical."}
{"text": "1. Field\nThe present disclosure relates generally to the electrical arts, and more particularly, to concepts and techniques for measuring the voltage of a power source.\n2. Background\nDetermining the voltage of a single cell battery is a straightforward procedure. A conventional voltmeter is simply placed across the terminals of the cell and the voltage measured. This procedure, however, poses various technological challenges when measuring the cell voltages of a multiple cell battery. In particular, the voltmeter must be switched between the cells to determine the voltage of each cell. Moreover, the voltmeter, which is generally composed of relatively low voltage breakdown semiconductor based electronic components, must withstand the voltage measured at each cell in the battery with respect to ground. This voltage, which is often referred to as “common mode voltage,” can reach hundreds of volts in large series connected battery stacks, such as those found in automobiles and other high voltage applications. These high voltage applications are beyond the voltage breakdown capabilities of most semiconductor components. Semiconductor based switches suffer from similar problems due to voltage breakdown limitations. Accordingly, there is a need in the art for isolated measurement techniques for batteries and other power sources."}
{"text": "This application claims the priority of German Patent Application No. 101 54 669.6, filed Nov. 7, 2001, the disclosure of which is expressly incorporated by reference herein.\nThe present invention relates to an internal combustion engine having at least two cylinder banks and more particularly, to an internal combustion engine whose cylinder heads are sealed by cylinder head covers, wherein to ventilate the crankcase from the so-called blow-by gases, ventilation lines are connected to the cylinder head covers and communicate with a negative pressure source, e.g., an intake pipe, and on the inside of the cylinder head cover means are provided for pre-separating the oil from the blow-by gases.\nU.S. Pat. No. 3,908,617 discloses a device for crankcase ventilation of an internal combustion engine with two cylinder banks in which ventilation lines mounted above the cylinder head housing or the cylinder head cover remove the blow-by gases located in the crankcase volume and return them to the intake system of the internal combustion engine in a closed circuit. In addition, sheet metal guide elements are mounted on the inside of the cylinder head cover. The blow-by gases flow past these guide elements and a portion of the oil carried along by the blow-by gases is deposited thereon."}
{"text": "The present invention relates to temperature-producing conductive-resistive medium and to a method of producing a variety of articles therefrom.\nThere have been many attempts to produce electrically-conductive coatings such as paints. Generally, there are two types of electrically-conductive coatings. The first is a low resistivity, high conductivity paint that contains a pigmentation of metal particles while the second is a high resistivity, low conductivity paint that is formed from compositions containing carbon or graphite.\nLow resistivity paints have traditionally been used to provide coatings having a high conductivity for connecting conductors that require a superior electrical bond with a minimum resistance. Generally, low resistivity paints cannot be applied to materials in order to produce temperature adjustable heating elements because the low resistivity paint requires a high volume of current to generate a reasonable output of heat. In contrast, the resistivity of traditional highly resistive paints is often so high that a relatively high voltage drop is required in order to generate sufficient heat. As a result, the use of high resistivity paints usually sacrifices safety. Furthermore, when either of the above-identified traditional conductive paints are applied to various substrates, cracks and flaking of the paint often develop over a period of time. This causes a breakdown in the temperature adjustable property of the article.\nIt is therefore an object of the present invention to provide a method and apparatus for generating an electrical resistance temperature adjustable substance for application to a variety of substrates in order to provide temperature controllable properties.\nIt is another object of the present invention to provide a method and apparatus for generating an electrical resistance temperature adjustable substance for application to a variety of materials wherein the electrical resistance temperature adjustable substance does not inhibit the inherent flexibility of the substrate to which it is applied.\nOther and further objects will be made known to the artisan as a result of the present disclosure and it is intended to include all such objects which are realized as a result of the disclosed invention."}
{"text": "The present invention relates generally to the field of solid state electrochemical devices, and more particularly to substrate, electrode and cell structures for solid state electrochemical devices.\nSolid state electrochemical devices are often implemented as cells including two porous electrodes, the anode and the cathode, and a dense solid electrolyte and/or membrane which separates the electrodes. For the purposes of this application, unless otherwise explicit or clear from the context in which it is used, the term “electrolyte” should be understood to include solid oxide membranes used in electrochemical devices, whether or not potential is applied or developed across them during operation of the device. In many implementations, such as in fuel cells and oxygen and syn gas generators, the solid membrane is an electrolyte composed of a material capable of conducting ionic species, such as oxygen ions, or hydrogen ions, yet has a low electronic conductivity. In other implementations, such as gas separation devices, the solid membrane is composed of a mixed ionic electronic conducting material (“MIEC”). In each case, the electrolyte/membrane must be dense and pinhole free (“gas-tight”) to prevent mixing of the electrochemical reactants. In all of these devices a lower total internal resistance of the cell improves performance.\nThe ceramic materials used in conventional solid state electrochemical device implementations can be expensive to manufacture, difficult to maintain (due to their brittleness) and have inherently high electrical resistance. The resistance may be reduced by operating the devices at high temperatures, typically in excess of 900° C. However, such high temperature operation has significant drawbacks with regard to the device maintenance and the materials available for incorporation into a device, particularly in the oxidizing environment of an oxygen electrode, for example.\nThe preparation of solid state electrochemical cells is well known. For example, a typical solid oxide fuel cell (SOFC) is composed of a dense electrolyte membrane of a ceramic oxygen ion conductor, a porous anode layer of a ceramic, a metal or, most commonly, a ceramic-metal composite (“cermet”), in contact with the electrolyte membrane on the fuel side of the cell, and a porous cathode layer of a mixed ionically/electronically-conductive (MIEC) metal oxide on the oxidant side of the cell. Electricity is generated through the electrochemical reaction between a fuel (typically hydrogen produced from reformed methane) and an oxidant (typically air). This net electrochemical reaction involves charge transfer steps that occur at the interface between the ionically-conductive electrolyte membrane, the electronically-conductive electrode and the vapor phase (fuel or oxygen). The contributions of charge transfer step, mass transfer (gas diffusion in porous electrode), and ohmic losses due to electronic and ionic current flow to the total internal resistance of a solid oxide fuel cell device can be significant. Moreover, in typical device designs, a plurality of cells are stacked together and connected by one or more interconnects. Resistive loss attributable to these interconnects can also be significant.\nIn work reported by de Souza and Visco (de Souza, S.; Visco, S. J.; De Jonghe, L. C. Reduced-temperature solid oxide fuel cell based on YSZ thin-film electrolyte. Journal of the Electrochemical Society, vol. 144, (no. 3), Electrochem. Soc, March 1997. p.L35-7. 7), a thin film of yttria stabilized zirconia (YSZ) is deposited onto a porous cermet electrode substrate and the green assembly is co-fired to yield a dense YSZ film on a porous cermet electrode. A thin cathode is then deposited onto the bilayer, fired, and the assembly is tested as an SOFC with good results. In work reported by Minh (Minh, N. Q. (Edited by: Dokiya, M.; Yamamoto, O.; Tagawa, H.; Singhal, S. C.) Development of thin-film solid oxide fuel cells for power generation applications. Proceedings of the Fourth International Symposium on Solid Oxide Fuel Cells (SOFC-IV), (Proceedings of the Fourth International Symposium on Solid Oxide Fuel Cells (SOFC-IV), Proceedings of Fourth International Symposium Solid Oxide Fuel Cells, Yokohama, Japan, 18-23 Jun. 1995.) Pennington, N.J., USA: Electrochem. Soc, 1995. p. 138-45), a similar thin-film SOFC is fabricated by tape calendaring techniques to yield a good performing device. However, these Ni—YSZ supported thin-film structures are mechanically weak, and will deteriorate if exposed to air on SOFC cool-down due to the oxidation of Ni to NiO in oxidizing environments. Also, nickel is a relatively expensive material, and to use a thick Ni—YSZ substrate as a mechanical support in a solid state electrochemical device will impose large cost penalties.\nSolid state electrochemical devices are becoming increasingly important for a variety of applications including energy generation, oxygen separation, hydrogen separation, coal gasification, and selective oxidation of hydrocarbons. These devices are typically based on electrochemical cells with ceramic electrodes and electrolytes and have two basic designs: tubular and planar. Tubular designs have traditionally been more easily implemented than planar designs, and thus have been preferred for commercial applications. However, tubular designs provide less power density than planar designs due to their inherently relatively long current path that results in substantial resistive power loss. Planar designs are theoretically more efficient than tubular designs, but are generally recognized as having significant safety and reliability issues due to the complexity of sealing and manifolding a planar stack.\nThus, solid state electrochemical devices incorporating current implementations of these cell designs are expensive to manufacture and may suffer from safety, reliability, and/or efficiency drawbacks. Some recent attempts have been made to develop SOFCs capable of operating efficiently at lower temperatures and using less expensive materials and production techniques. Plasma spray deposition of molten electrolyte material on porous device substrates has been proposed, however these plasma sprayed layers are still sufficiently thick (reportedly 30-50 microns) to substantially impact electrolyte conductance and therefore device operating temperature.\nAccordingly, a way of reducing the materials and manufacturing costs and increasing the reliability of solid state electrochemical devices would be of great benefit and, for example, might allow for the commercialization of such devices previously too expensive, inefficient or unreliable."}
{"text": "Modern farmers strive to improve the management of increasing amounts of farm acres. Improving management requires farmers to be able to quickly prepare the soil for each season's farming operations. This haste has driven the need for more efficient and larger farming equipment.\nImplements such as harrows, packers, or combined harrow-packers were some of the earliest implements to be made with widths exceeding sixty feet in the field operating position. As tractor horsepower has increased over time, larger tillage implements have been made available. These larger implements require a mechanism for compactly folding the implement for practical and safe transport over the highway. U.S. Pat. No. 4,821,809, patented by Summach et al., discloses a convenient mechanism for such folding.\nThe conventional method of folding tillage implements is by folding wing sections along forward aligned axes such that the wings are folded to a generally upright position. Double folding wing sections may have outer sections that fold inwardly and downwardly from the ends of inner wing sections in five section winged implements. In the case of these conventional wing implements, the minimum implement width that can be achieved by such folding is limited by the width of the center section. As a result, road transport may still be somewhat restricted as these implements often exceed twenty feet or more in transport width.\nRoad transport standards in North America are beginning to follow the standards set in Europe in which maximum road transport widths and heights for agricultural implements are being defined. Large implements that have conventional folding wing sections are not able to be folded such that they fall within width and height limits that may be generally 3 meters wide and 4 meters high. Some U.S. states have adopted transport width limits of 13.5 ft.\nForward or rear folding implements provide some relief with respect to such transport limits. However, implements must also be made to function with the accurate seeding ability that conventionally folded implements have become capable of. Although some rear or forward folding multibar tillage implements have been developed, they do not demonstrate the accurate depth control required for farming operations.\nOne problem is that a tillage-packer combination for drill seeding requires the gang supporting tillage elements to be maintained parallel to the ground through a range of adjustable operating levels. The drawbar disclosed in Summach '809 raises or lowers the first attached gang of elements in a rotatable manner through its field and transport ranges of motion. A level manner of height adjustment is required for tillage elements.\nAnother problem that must be overcome for compact folding is the avoidance of the packer elements of the second gang striking the tillage elements of the first gang when raised to the transport position. If compact folding is not required, then the downward rotation of the suspended second gang may be limited so as not to impact the elements of the first gang. But when compact folding is desired, the elements of the second gang are in direct alignment with the ground elements of the first gang so that alignment is achieved.\nTherefore, a multibar implement is required for the tillage of high acre farms with both speed and efficiency."}
{"text": "The present invention relates to a new and distinct cultivar of Abelia of hybrid origin and will be referred to hereafter by its cultivar name, ‘ABENOV41’. ‘ABENOV41’ represents a new Abelia, a shrub grown for landscape use.\nThe new cultivar is the result of a controlled breeding program conducted by the Inventor in Angers-Beaucouzé, France. The objective of the breeding program is to develop new cultivars of Abelia twith increased cold hardiness and unique flower colors and plant habits.\nThe new cultivar arose from a cross made by the Inventor in June of 2004 between Abelia schumannii ‘Bumblebee’ (not patented) as the female parent and Abelia×grandiflora ‘Semperflorus’ (not patented) as the male parent. The Inventor selected ‘ABENOV41’ as a single unique plant amongst the seedlings that resulted from the above cross in 2006.\nAsexual propagation of the new cultivar was first accomplished by softwood stem cuttings by the Inventor in Angers-Beaucouzé, France in September of 2006. Asexual propagation by softwood stem cuttings has determined that the characteristics of the new cultivar are stable and are reproduced true to type in successive generations."}
{"text": "Photocatalysts employing titanium oxide are widely used by virtue of low price; high chemical stability, high photocatalytic activity (organic compound degradability, anti-bacterial property, etc.); non-toxicity to the human body; etc.\nIt has been known that a mixture of titanium oxide with metallic copper or a copper compound, or a product of titanium oxide on which copper or a copper compound has been deposited, serves as an excellent photocatalyst or an excellent anti-viral agent.\nFor example, Patent Document 1 discloses use of nano particles of a compound MnXy for suppression and/or prevention of infection with viruses. Examples of the nano-particle compound include TiO2, Cu2O, CuO, and combinations thereof.\nRegarding the aforementioned combinations of titanium oxide with metallic copper or a copper compound, particularly, the anti-viral performance of the photocatalysts has been enhanced by use of anatase-type titanium oxide.\nFor example, Patent Document 2 discloses an anti-bacterial photocatalytic aqueous coating material in which a metal such as copper is deposited on a photocatalyst such as titanium oxide. Patent Document 2 also discloses that titanium oxide preferably has an anatase-type crystal structure.\nPatent Document 3 discloses a phage/virus inactivating agent formed of anatase-type titanium oxide containing copper at a CuO/TiO2 (ratio by mass %) of 1.0 to 3.5. The invention of Patent Document 3 was accomplished with respect to the finding that copper-containing anatase-type titanium oxide can inactivate phages/viruses.\nRegarding the aforementioned combinations of titanium oxide with metallic copper or a copper compound, it have been known that a monovalent copper compound has particularly used as a copper compound, and the monovalent copper exhibits excellent microorganism- and virus-inactivating performance.\nFor example, Patent Document 4 discloses an anti-viral coating material, characterized by containing a monovalent copper compound as an active ingredient which can inactivate viruses. Patent Document 4 also discloses that the monovalent copper compound inactivates a variety of viruses through contact therewith.\nPatent Document 5 discloses a microorganism inactivating agent which contains a monovalent copper compound as an active ingredient for use in inactivation of microorganisms in a short time. Patent Document 5 also discloses another microorganism inactivating agent which contains a monovalent copper compound and a photocatalytic substance. The photocatalytic substance may be a titanium oxide catalyst. Patent Document 5 further discloses that a monovalent copper compound exhibits a remarkably strong microorganism-inactivating effect, as compared with a divalent copper compound.    Patent Document 1: Japanese Kohyo (PCT) Patent Publication No. 2009-526828    Patent Document 2: Japanese Patent Application Laid-Open (kokai) No. 2000-95976    Patent Document 3: Japanese Patent No. 4646210    Patent Document 4: Japanese Patent Application Laid-Open (kokai) No. 2010-168578    Patent Document 5: Japanese Patent Application Laid-Open (kokai) No. 2011-190192"}
{"text": "1. Field of the Invention\nThe invention relates to a method for multiple exposure at least of one substrate coated with a photosensitive layer, a microlithography projection exposure installation for multiple exposure at least of one substrate coated with a photosensitive layer, and a projection system having an illumination system and a projection objective.\n2. Description of the Related Art\nThe efficiency of projection exposure installations for the microlithographic fabrication of semiconductor components and other finely patterned devices is substantially determined by the imaging properties of the projection objectives. Moreover, the image quality and the wafer throughput that can be achieved with the installation are influenced substantially by properties of the illumination system arranged upstream of the projection objective. Said system must be capable of preparing the light of a primary light source, for example a laser, with the highest possible efficiency, and in the process of producing as uniform as possible a distribution of intensity in an illumination field of the illumination system.\nDepending on the nature and size of the patterns to be produced on the wafer, suitable exposure parameters can be set on the illumination system and/or the projection objective. For example, conventional illumination with different degrees of coherence and annular field illumination or polar illumination can be set on the illumination system in order to produce an off-axis, oblique illumination. The numerical aperture can be set on the projection objective.\nGiven a prescribed wavelength of the primary light source, the selection of suitable exposure parameters can serve, inter alia, for imaging structures that it would not be possible to image with satisfactory quality by using other exposure parameters because of their small structural sizes. However, a long exposure time is frequently associated with a selection of exposure parameters for which such fine structures can be resolved, and so the wafer throughput turns out to be low. Such fine structures frequently cannot be correctly resolved given a selection of exposure parameters for which a higher wafer throughput is achieved, although structures with larger structural sizes can be.\nSince the structures to be produced on the wafer can frequently be subdivided into fine and coarse structures, it can be favorable to carry out a double exposure of the wafer for which a first set of exposure parameters is used for imaging the coarse structures, while a second set of exposure parameters different from the first one is used for imaging the fine structures. The exposure parameters of the first set can be selected, for example, such that only a short exposure time is required for imaging the coarse structures. The exposure parameters of the second set can be optimized such that only those structures are imaged that are so fine that they cannot be imaged with the first set of exposure parameters. Of course, multiple exposures with more than two exposures are also possible.\nIn a known type of double exposure, an exposure by means of an amplitude mask is carried out with a first set of exposure parameters. It is possible in this way to make use of oblique illumination such as, for example, annular, dipole or quadrupole illumination in order to increase the resolution. A second exposure with a second set of exposure parameters is carried out with the aid of a phase mask. A coherent illumination with a low degree of coherence σ is normally set thereby at the illumination system. Such a method is described, for example, in the article entitled “Improving Resolution in Photolithography with a Phase-Shifting Mask” by M. D. Levenson, N. S. Viswanathan, R. A. Simson in IEEE Trans. Electr. Dev., ED-29(12), pp. 1828-1836, 1982, and in the article entitled “Performance Optimization of the Double-Exposure” by G. N. Vandenberghe, F. Driessen, P. J. van Adrichem, K. G. Ronse, J. Li, L. Karlaklin in Proc. of the SPIE, Vol. 4562, pp. 394-405, 2002. In the method for multiple exposure that is described there, a first set of exposure parameters is initially set on the projection system for the first exposure. Thereafter, a second set of exposure parameters is set on the projection system for the second exposure, and this requires a reconfiguration of the projection system. This reconfiguration of the projection system from the first to the second set of exposure parameters results in a time loss and in mechanical wear of the parts whose position and/or shape need to be varied in the reconfiguration."}
{"text": "From DE 1 103 216 a device for distributing cut tobacco to cigarette-making machines is known, wherein the cut tobacco is fed from a conveyor onto a rotary table from which the tobacco is drawn by stationary sucking pipes spaced at the periphery of a table constituting a distributing element, the cut tobacco fed from the conveyor falling onto a cone located centrally relative to the rotary table. The cut tobacco slides down along the cone onto the rotary table gravitationally and then it is transported due to the centrifugal force as a layer towards the periphery of the table, from where it is sucked by vertical pipes to deliver the cut tobacco to the cigarette-making machines.\nDE 198 23 873 presents a similarly operating device for feeding cut tobacco to many machines. The cut tobacco is fed via a vertical channel onto a bowl performing a composed, rotary and circulating, motion. The sucking channels, picking up the cut tobacco from the uniformly formed layer, are arranged vertically within the bowl cover at the bowl periphery.\nIn GB 959 343 a device is described in which the cut tobacco is fed, as previously, from above onto a rotary distribution disk and is directed by the centrifugal force towards receiving channels arranged radially in the side wall of the distribution chamber.\nIn a slightly different arrangement, known from DE 300 90 000, cut tobacco is fed through a charging hopper onto a linear vibrational conveyor. The vibrational conveyor transfers the fed cut tobacco to a place above which sucking pipes are situated. The cut tobacco is transported in the form of a layer and the sucking pipes are arranged vertically just above the surface of this layer.\nUsually the bottom of the distribution chamber is flat or has the shape of a bowl and it is a surface of revolution and posses a centrally located rotational cone.\nThe process of feeding the cut tobacco to the cigarette-making machines is discontinuous, the result of which is that the more receiving channels are connected, the more frequent changes of the flow rate of the tobacco through the distributing device will occur. The discontinuity of the feeding process results from the fact that after filling the cut tobacco container located within the machine, the feeding is stopped until the amount of the cut tobacco in the container drops below a certain predefined level, afterwards the feeding is started again. Devices for distributing cut tobacco, employed in the tobacco industry, usually feed a lot of cigarette-making machines. Every change in a total throughput of the receiving channels will result, as a consequence, in a change of the efficiency of the conveyor feeding the distributing device.\nAll the solutions presented above relate to devices for distributing cut tobacco to cigarette-making machines using gravitational feeding, usually in the form of a feeding channel and a couple of pneumatic receiving channels transferring the cut tobacco to the cigarette-making machines, the receiving channels being connected to the distributing chamber or being located at the periphery of the distributing element for uniform distributing the cut tobacco into the inlets of the receiving channels. For proper operation of all the above devices it is necessary to collect some amount of the cut tobacco in the distribution chamber, which is transferred to the space from which it is received by the receiving channels. During transferring the layer of the cut tobacco gains its optimal thickness in order to ensure repeatable conditions of receiving the cut tobacco by the receiving channels. Therefore the receiving channels are distant from the feeding channel. In each of the devices in the case of temporary stopping the process of feeding the cigarette-making machines, the amount of the cut tobacco, which has been already delivered to the distributing device but has not been yet received, is an excess of the cut tobacco present in the device relative to the amount necessary for its operation. The cut tobacco tends to agglomerate, i.e., to create bundles, the effect of the agglomeration being particularly strong if the cut tobacco is stored in a high layer, as in the vertical channel feeding the distributing device.\nIf the process of receiving the cut tobacco by the cigarette-making machines, connected to a single distributing device, is stopped, one must stop the conveyor feeding the device, which was operating with a rate adjusted for feeding all the cigarette-making machines. However, due to inertia of the system, the distribution chamber will be filled anyway as well as, partially or fully, then vertical feeding channel. Restarting the device after a longer downtime may occur difficult, since the bulk density of the cut tobacco collected and stored under a pressure within the feeding channel increases and it is significantly more difficult to form a uniform layer of the cut tobacco and to suck the agglomerated tobacco through the receiving channels. Sometimes, in order to restart the feeding system the agglomerated tobacco must be removed from the lower portion of the feeding channel and partially from the distribution chamber.\nIf a couple of receiving channels will be shut off simultaneously, i.e., in the case of a rapid drop of the received amount of the cut tobacco, an excess of the cut tobacco will arise within the distribution chamber. The efficiency of the conveyor feeding the distributing device will be adjusted to the throughput of the cigarette-making machines that are still working, and the excess of the collected cut tobacco will be used by those machines, however if the excess is relatively large, disturbances in the receiving process may arise.\nFrequently, cigarette manufacturers must face the task of producing short series of new cigarette brands. Large distributing devices with rotary tables or vibrational conveyors are expensive and there is no economical justification for using them in the case of frequent changes of the brand of tobacco fed to one or two cigarette-making machines."}
{"text": "Die stacking integrates semiconductor devices vertically in a single package in order to directly influence the amount of silicon that can be included in a given package footprint. Die stacking simplifies the surface-mount pc-board assembly and conserves pc-board real estate because fewer components are placed on the board. Die stacking has included different memory combinations that place flash memory with SRAM and RAM. Die stacking has evolved to multiple die stacks and side-by-side combinations of stacked and unstacked dies within a package.\nThe dies are mounted on a substrate which may then be bumped to create either a Chip Scale Package (CSP) or a Ball Grid Array (BGA) as the final package. Present die stacking techniques include mounting smaller dies onto larger ones to enable wire bonding of both, as well as techniques for stacking same-size die. To further increase the memory density and memory bandwidth available in a given size footprint, Package-on-Package (PoP) may be utilized to vertically connect multiple packages such as a logic package with a memory package, where each package may contain one or more die. Still, improvements in packaging are needed as the number of stacked-die in a package is expanded.\nIt will be appreciated that for simplicity and clarity of illustration, elements illustrated in the figures have not necessarily been drawn to scale. For example, the dimensions of some of the elements may be exaggerated relative to other elements for clarity. Further, where considered appropriate, reference numerals have been repeated among the figures to indicate corresponding or analogous elements."}
{"text": "The present invention relates to shaft coupling and is particularly concerned with the temporary coupling of shafts using fluid coupling. It finds particular applications in gas turbine engines where shafts may be rotating at very different speeds but may equally find utility in other applications.\nIt is commonplace in gas turbine engines used for aircraft propulsion to remove power from a shaft of an engine during flight to drive aircraft systems and accessories including cabin systems, in-flight entertainment systems and cabin air pressurisation systems. In a three-shaft gas turbine engine it is known to take power from the high pressure (HP) shaft. Since more power is usually available from the intermediate pressure (IP) shaft than the HP shaft, it is beneficial to take power from the IP shaft, particularly during engine descent, and hence reduce fuel consumption compared to the condition if power is taken from the HP shaft. There is a consequent reduction in the amount of fuel required for a flight and thus the cost of that flight. However, there is a requirement to drive the high pressure compressor during engine starting meaning it may be desirable to transfer starting torque from the IP to the HP shaft. Furthermore, under high power off-take conditions when the engine is at idle it may be necessary to transfer power from the HP to the IP shaft. Therefore, there may be a need to couple the high pressure and intermediate pressure shafts for some periods in the engine cycle.\nA conventional method of temporarily coupling two shafts is to use fluid coupling. A typical arrangement is shown in FIG. 1 in which a first, intermediate pressure shaft 10 and a second, high pressure shaft 12 are coupled by fluid couplings 14. Typically the fluid couplings 14 comprise two or more ball chambers annularly arrayed around the ends of shafts 10, 12. The chambers 14 are linked by fluid passages 16 so that the working fluid, for example oil, can be distributed into and between the chambers 14 or removed from them. Each chamber 14 and the associated passages 16 is formed in two sections, connected to the IP and HP shafts 10, 12 respectively. When the chambers 14 and passages 16 are filled with oil the two sections of the chambers 14a, 14b are compelled to rotate in approximate synchronicity and thus the shafts 10, 12 are coupled. When decoupling is required, the oil is drained from the chambers 14 and passages 16 so that the two sections are no longer constrained to rotate together but are free to rotate in synchronicity with their respective shafts, which may be rotating at different speeds. Hence the two sections of each chamber 14 rotate independently at different speeds to each other."}
{"text": "1. Field of the Invention\nThe present invention generally relates to methods and systems for inspection of an entire wafer surface using multiple channels. Certain embodiments relate to detecting light scattered from different portions of the entire wafer surface using different detection channels.\n2. Description of the Related Art\nFabricating semiconductor devices such as logic and memory devices typically includes processing a specimen such as a semiconductor wafer using a number of semiconductor fabrication processes to form various features and multiple levels of the semiconductor devices. For example, lithography is a semiconductor fabrication process that typically involves transferring a pattern to a resist arranged on a semiconductor wafer. Additional examples of semiconductor fabrication processes include, but are not limited to, chemical-mechanical polishing, etch, deposition, and ion implantation. Multiple semiconductor devices may be fabricated in an arrangement on a semiconductor wafer and then separated into individual semiconductor devices.\nWafers may contain defects both in central portions of the wafers as well as in edge portions of the wafers, which includes a relatively narrow region around the periphery of the wafers, and on the outer edge of the wafers. Examples of defects that may be found in the edge portion and on the outer edge of wafers include, but are not limited to, chips, cracks, scratches, marks, particles, and residual chemicals (e.g., resist and slurry). As wafer sizes continue to increase, both wafer and integrated circuit (IC) manufacturers are becoming more concerned about defectivity at or near the wafer edge. The main concerns are that edge defects could fall onto the central part of the wafer thereby causing untraceable yield loss, cross contamination during processing, and/or catastrophic wafer breakage. These yield loss mechanisms are experienced by most wafer and IC manufacturers at one time or another.\nTraditionally, wafer inspection tools are designed to inspect a central portion of the wafers (i.e., a surface area of the wafer on which electrical elements will be formed or a surface area of the wafer opposite that on which electrical elements will be formed). Since these areas of the wafer reflect or scatter relatively small amounts of light, such wafer inspection tools are designed to detect relatively small amounts of light. However, near the outer edge of the wafer, relatively large amounts of light may be reflected or scattered from the wafer due to edge features such as a bevel formed at or near the outer edge. As a result, these large amounts of light will saturate the detectors of traditional wafer inspection systems. Consequently, any output signals generated near or at the edge of wafers by such wafer inspection tools are generally unusable. In some instances, the wafer inspection systems may be designed to block the light from reaching the detectors when inspecting near the edge of the wafer to protect the detectors from damage that may be caused by the relatively high intensity light.\nSome edge inspection systems are being developed to detect defects at or near the outer edge of wafers. Examples of apparatuses for detecting defects along the edge of electronic media such as semiconductor wafers are illustrated in U.S. Patent Application Publication Nos. 2003/0030050 by Choi and 2003/0030795 by Swan et al., which are incorporated by reference as if fully set forth herein. Due to the substantially different reflecting and scattering characteristics of the outer edge of wafers in comparison to the inner portion of the wafer, such edge inspection systems have substantially different configurations than the traditional wafer inspection tools. Therefore, the edge inspection systems are not optimized to, or even able to, detect defects in the central portion of the wafers. Consequently, if wafer or IC manufacturers want to detect defects in both the central and outer portions of wafer (as is usually the case since defects in either portion may result in expensive yield losses and other problems), they will need to purchase two separate tools. Using two different wafer inspection tools instead of just one inspection tool will obviously increase costs in many ways such as increases in clean room real estate and operating costs, increases in tool maintenance costs, and increases due to reduced throughput. However, since a tool that is capable of inspecting both the inner and outer portions of wafers is not currently available, and due to the increasing costs associated with defect-based yield losses, wafer and IC manufacturers may not be able to avoid the costs associated with multiple, different inspection tools.\nAccordingly, it may be advantageous to develop a wafer inspection system that is capable of inspecting substantially an entire surface of wafers including both center and edge portions of the wafers."}
{"text": "This invention relates generally to a method and system for facilitating the identification, investigation, assessment and management of legal, regulatory financial and reputational risks (“risks”). In particular, the present invention relates to a computerized system and method for banks and non-bank financial institutions to manage risks associated with maintaining investment accounts for a politically identified person (PIP).\nRisk associated with maintaining an investment account can include factors associated with financial risk, legal risk, regulatory risk and reputational risk. Financial risk includes factors indicative of monetary costs that the financial institution may be exposed to as a result of opening a particular account and/or transacting business with a particular client. Monetary costs can be related to fines, forfeitures, cost to defend an adverse position, or other related potential sources of expense. Regulatory risk includes factors that may cause the financial institution to be in violation of rules put forth by a regulatory agency such as the Securities and Exchange Commission (SEC). Reputational risk relates to harm that a financial institution may suffer regarding its professional standing in the industry. A financial institution can suffer from being associated with a situation that may be interpreted as contrary to an image of honesty and forthrightness.\nRisk associated with an account for a PIP can be greatly increased due to the nature of a position held by the PIP along with the power and knowledge associated with that position. PIPs can include an elected official, a bureaucrat, a political appointee, a World Bank Official, a military person, or other individual associated with a sovereign power or international organization. In addition, a PIP can be a person who holds a position in the private sector wherein the position is associated with politically sensitive influences. As part of due diligence associated with managing financial accounts, it is imperative for a financial institution to “Know their Customer” including such attributes as a position held by the customer and the magnitude of risk associated with that position.\nCompliance officers and other financial institution personnel typically have few resources available to assist them with the identification of present or potential risks associated with a particular investment or trading account associated with a PIP. Risk can be multifaceted and far reaching. The amount of information that needs to be considered to evaluate whether an individual is a PIP, and whether a particular PIP poses a significant risk, is substantial, if not overwhelming. A multi-agency working group drawn from the U.S. Justice and Treasury Departments and various federal agencies issued guidelines to meet the perceived threat arising from the transmission of the proceeds of foreign political corruption to U.S. financial institutions. These guidelines are designed to counsel banks, broker-dealers, and other financial institutions on their obligations with regard to funds that appear to be related to the theft of sovereign assets by foreign political leaders. The guidelines follow in the wake of embarrassing disclosures relating to the transmission of funds to U.S. financial institutions by a long list of foreign leaders including the Marcos family of the Philippines, Raul Salinas of Mexico, and General Abacha of Nigeria.\nHowever, financial institutions do not have available a mechanism which can provide real time assistance to assess a risk factor associated with a PIP, or otherwise qualitatively manage such risk. In the event of investment problems, it is often difficult to quantify to regulatory bodies, shareholders, newspapers and/or other interested parties, the diligence exercised by the financial institution to properly identify and respond to risk factors. Absent a means to quantify good business practices and diligent efforts to contain risk, a financial institution may appear to be negligent in some respect.\nWhile the guidelines offered by the U.S. federal government are only advisory and therefore not a law, rule or regulation, according to the guidelines, deficiencies in a financial institution's anti-money laundering controls may prompt a regulator to require that the guidance be integrated into the institutions policies and procedures.\nWhat is needed is a method and system to assist with risk management and due diligence related to financial accounts associated with a PIP. A new method and system should anticipate offering guidance to personnel who interact with clients and help the personnel identify PIP. In addition, it should be situated to convey risk information to a compliance department and be able to demonstrate to regulators that a financial institution has met standards relating to risk containment."}
{"text": "Glucose-regulating peptides are critical regulatory components of human metabolism. Various peptides have been described with biological effects that result in either an increase or decrease in serum glucose levels. These peptides tend to be highly homologous to each other, even when they possess opposite biological functions. Many glucose regulating peptides, including those used as therapeutics, are typically labile molecules exhibiting short shelf-lives, particularly when formulated in aqueous solutions. In addition, many glucose regulating peptides have limited solubility, or become aggregated during recombinant productions, requiring complex solubilization and refolding procedures. Various chemical polymers can be attached to such peptides and proteins to modify their properties. Of particular interest are hydrophilic polymers that have flexible conformations and are well hydrated in aqueous solutions. A frequently used polymer is polyethylene glycol (PEG). These polymers tend to have large hydrodynamic radii relative to their molecular weight (Kubetzko, S., et al. (2005) Mol Pharmacol, 68: 1439-54), and can result in enhanced pharmacokinetic properties. However, the chemical conjugation of polymers to proteins requires complex multi-step processes; typically, the protein component needs to be produced and purified prior to the chemical conjugation step and the conjugation step can result in the formation of heterogeneous product mixtures that need to be separated, leading to significant product loss. Alternatively, such mixtures can be used as the final pharmaceutical product, but are difficult to standardize. Some examples are currently marketed PEGylated Interferon-alpha products that are used as mixtures (Wang, B. L., et al. (1998) J Submicrosc Cytol Pathol, 30: 503-9; Dhalluin, C., et al. (2005) Bioconjug Chem, 16: 504-17). Such mixtures are difficult to reproducibly manufacture and characterize as they contain isomers with reduced or no therapeutic activity.\nAlbumin and immunoglobulin fragments such as Fc regions have been used to conjugate other biologically active proteins, with unpredictable outcomes with respect to increases in half-life or immunogenicity. Unfortunately, the Fc domain does not fold efficiently during recombinant expression and tends to form insoluble precipitates known as inclusion bodies. These inclusion bodies must be solubilized and functional protein must be renatured. This is a time-consuming, inefficient, and expensive process that requires additional manufacturing steps and often complex purification procedures.\nThus, there remains a significant need for compositions and methods that would improve the biological, pharmacological, safety, and/or pharmaceutical properties of glucose regulating peptides."}
{"text": "It is known to fill containers while they are arranged in a space that is separated from the environment by a housing. It is also known to expose the containers to laminar flow of sterile air directed from the top down while the containers are in this space. Doing so helps to avoid contamination of the containers with foreign particles, such as dust, germs, or other unwanted substances.\nIt is also known to have transport sections between container handling machines or assemblies. These container-handling machines convey containers from one container-handling machine to a subsequent container-handling machine along a container-transport direction. A disadvantage of known transport sections is that any empty containers, or containers that have bee filled but not yet sealed, remain un protected against contamination while being transported along a transport section."}
{"text": "1. Field of the Invention\nThe present invention relates to a channel allocation method of a wireless network and a system thereof. More particularly, the present invention relates to a distributed channel allocation method of a wireless mesh network and a system thereof.\n2. Description of Related Art\nIn recent years, there is a rapid development in the field of wireless broadband access techniques including Wi-Fi (IEEE 802.11 series), WiMAX (IEEE 802.16 series) and 3G, etc. The wireless mesh network (referred to hereinafter as WMN, IEEE 802.11s) is one of the key techniques integrated with the wireless broadband network. The structure of the WMN illustrated in FIG. 1 is a mesh network based on a wireless transmission interface, and the WMN has a similar operation mode to that of an Ad-hoc network. Since the operation of the WMN is based on the wireless transmission interface, it has the advantage of rapid deployment without restriction of the geographical landforms. The WMN is generally applied to a community network, a temporary network of exhibition halls or shopping stalls, networks established in disaster areas or areas having special geographical environments, and so on.\nThe operation of the WMN is based on the wireless transmission interface. Taking the IEEE 802.11a/g for an example, its transmission bandwidth of data is 54 Mbps (mega bytes per second), which is the maximum possible transmission bandwidth. However, influenced by a MAC (media access control) contention, 802.11 headers, 802.11 ACK signals and packet errors, an average applicable bandwidth is usually less than half of the maximum bandwidth.\nFurthermore, the most serious issue lies in that a data transmission rate of a network link layer may be decreased greatly due to signal interference. Two possible interference problems are shown in FIG. 2: (1) interference in the same transmission path, (2) interference in the adjacent transmission paths. Referring to FIG. 2, the signal coverage of a node 3 includes nodes 2, 4 and 9. Similarly, the node 3 is simultaneously in the signal coverage of the nodes 2, 4 and 9. A first transmission path and a second transmission path are paths for data transmission. The first transmission path is taken for an example. When the node 2 and the node 3 are transmitting data, the node 4 may receive signals from the node 3, resulting in the fact that node 4 cannot transmit data to a node 5 provisionally. Therefore, the bandwidth of the first transmission path is reduced, which refers to the so-called interference in the same transmission path.\nOn the other hand, referring to the node 9 on the second transmission path, since the node 9 is in the signal coverage of the node 3, the node 9 may receive signals from the node 3 when the node 2 and the node 3 are transmitting data, resulting in the fact that the node 9 cannot transmit data to a node 8 or a node 10 provisionally. The phenomenon indicating an interference of data transmission through the first transmission path with that through another transmission path (a second transmission path) represents the so-called interference in the adjacent transmission paths. Therefore, many studies are performed on the WMN to learn how to improve an applicable bandwidth of the WMN by advancing a structural design thereof.\nAccording to the IEEE 802.11s WiFi Mesh standard, a plurality of wireless transmission interfaces is allowed to use different non-overlapping channels for transmission, so as to increase the transmission bandwidth. Therefore, some studies have been developed to increase a network flow by applying multi-network interface cards (referred to hereinafter as Multi-NIC). A method of increasing the network flow includes allocating a plurality of NICs on each node, and each of the NICs may employ a different non-overlapping channel to communicate with other nodes. The advantage of this method lies in that it is unnecessary to modify any existing hardware structures. Only is an integral channel allocation method required for assisting the existing hardware structure, and the network flow can be substantially improved.\nA method and a system for assigning channels in a wireless local area network (WLAN) is disclosed in U.S. Publication No. 2006/0072502 A1, in which the WLAN infrastructure mode (i.e. a client to hub communication mode) is provided. A mobile node (referred to hereinafter as MN) in the network is connected to an access point (referred to hereinafter as AP) by means of one hop, and the other end of the AP is connected to a wired network, wherein each AP has at least two applicable channels, and each AP is at least adjacent to another AP.\nEach AP constantly collects the traffic load information and forecasts a possible throughput on each channel. Thereafter, the AP determines an optimal channel for connecting with the MN within the signal coverage of the AP. However, this channel allocation method only takes the optimal channel within one hop between the AP and the MN into account. Therefore, the application of the method is limited.\nMost of the early studies focus on modifying an MAC layer protocol of the network to support a multiple channel transmission. The studies aim to find the optimal channel for transmitting every single packet, so as to avoid the interference. On the other hand, a concept of a Multi-NIC disclosed by V. Bahl et al. and P. H. Hsiao et al. in two articles has drawn attention and discussions recently. One of the articles was authored by V. Bahl, A. Adya, J. Padhye, A. Wolman, entitled “Reconsidering the Wireless LAN Platform with Multiple Radios” Workshop on Future Directions in Network Architecture (FDNA-03), while another one was authored by P. H. Hsiao, A. Hwang, H. T. Kung, and D. Vlah, entitled “Load-Balancing Routing for Wireless Access Networks” Proc. of IEEE Infocom 2001. The method disclosed therein is to install a plurality of the NICs on each node of the Ad-hoc network, and each NIC may dynamically determine a channel for communicating with other nodes. The advantage of this method lies in that it is unnecessary to modify any existing hardware structures. Only is the integral channel allocation method required for assisting the existing hardware structure, and the network flow can be substantially improved. Sequentially, a channel allocation method based on a centralize structure was disclosed by A. Raniwala, K. Gopalan, T. Chiueh, entitled “Centralized channel assignment and routing algorithms for multi-channel wireless mesh networks,” ACM Mobile Computing and Communications Review 8 (2) (2003), which is one of the earliest articles having a formal definition of the channel allocation. In the method, a load-aware channel assignment is performed by an evaluation matrix defined by the authors themselves, the entire network is calculated in overall, and a preferable channel allocation is obtained. Thus, a maximum network flow is then achieved.\nIn recent studies, a channel allocation technique based on a dynamic & distributed structure has been disclosed, wherein channel allocation information is exchanged by using a common channel framework according to the IEEE 802.11s standard. This technique is based on IEEE 802.11 WLAN standard, wherein a plurality of wireless NICs is installed to support a multi-channel transmission. However, the interference still cannot be avoided in the aforementioned techniques."}
{"text": "When a service provider of communications deploys a communications network, there can be many challenges. Among them include without limitation the cost of deployment, the cost of adding network equipment to accommodate subscriber growth, maintenance of the network, serviceability of the network, and managing the addition of new subscribers to the network—just to mention a few."}
{"text": "Power amplifiers for cellular handsets are optimized for efficiency at, or close to, maximum output power. However, in the field, they may only be called upon to operate near maximum output power for a very small percentage of the time. The rest of the time, they may be operating at back-off output power levels, where their direct current (DC) to radio-frequency (RF) conversion efficiency is very much reduced. This reduced efficiency under practical conditions results in wasted battery power in the handset and, therefore, reduced talk time."}
{"text": "1. Field of the Invention\nThis application relates to the general field of Integrated Circuit (IC) devices and fabrication methods, and more particularly to multilayer or Three Dimensional Integrated Circuit (3D IC) devices and fabrication methods.\n2. Discussion of Background Art\nPerformance enhancements and cost reductions in generations of electronic device technology has generally been achieved by reducing the size of the device, resulting in an enhancement in device speed and a reduction in the area of the device, and hence, its cost. This may be generally referred to as ‘device scaling’. The dominant electronic device technology in use today may be the Metal-Oxide-Semiconductor field effect transistor (MOSFET) technology.\nPerformance and cost are driven by transistor scaling and the interconnection, or wiring, between those transistors. As the dimensions of the device elements have approached the nanometer scale, the interconnection wiring now dominates the performance, power, and density of integrated circuit devices as described in J. A. Davis, et. al., Proc. IEEE, vol 89, no. 3, pp. 305-324, March 2001 (Davis).\nDavis further teaches that three dimensional integrated circuits (3D ICs), i.e. electronic chips in which active layers of transistors are stacked one above the other, separated by insulating oxides and connected to each other by metal interconnect wires, may be the best way to continue Moore's Law, especially as device scaling slows, stops, or becomes too costly to continue. 3D integration would provide shorter interconnect wiring and hence improved performance, lower power consumption, and higher density devices.\nOne approach to a practical implementation of a 3D IC independently processes two fully interconnected integrated circuits including transistors and wiring, thins one of the wafers, bonds the two wafers together, and then makes electrical connections between the bonded wafers with Thru Silicon Vias (TSV) that may be fabricated prior to or after the bonding. This approach may be less than satisfactory as the density of TSVs may be limited, because they may require large landing pads for the TSVs to overcome the poor wafer to wafer alignment and to allow for the large (about one to ten micron) diameter of the TSVs as a result of the thickness of the wafers bonded together. Additionally, handling and processing thinned silicon wafers may be very difficult and prone to yield loss. Current prototypes of this approach only obtain TSV densities of 10,000 s per chip, in comparison to the millions of interconnections currently obtainable within a single chip.\nBy utilizing Silicon On Insulator (SOI) wafers and glass handle wafers, A. W. Topol, et. al., in the IEDM Tech Digest, p363-5 (2005), describe attaining TSVs of tenths of microns. The TSV density may be still limited as a result from misalignment issues resulting from pre-forming the random circuitry on both wafers prior to wafer bonding. In addition, SOI wafers are more costly than bulk silicon wafers.\nAnother approach may be to monolithically build transistors on top of a wafer of interconnected transistors. The utility of this approach may be limited by the requirement to maintain the reliability of the high performance lower layer interconnect metallization, such as, for example, aluminum and copper, and low-k intermetal dielectrics, and hence limits the allowable temperature exposure to below approximately 400° C. Some of the processing steps to create useful transistor elements may require temperatures above about 700° C., such as activating semiconductor doping or crystallization of a previously deposited amorphous material such as silicon to create a poly-crystalline silicon (polysilicon or poly) layer. It may be very difficult to achieve high performance transistors with only low temperature processing and without mono-crystalline silicon channels. However, this approach may be useful to construct memory devices where the transistor performance may not be critical.\nBakir and Meindl in the textbook “Integrated Interconnect Technologies for 3D Nanosystems”, Artech House, 2009, Chapter 13, illustrate a 3D stacked Dynamic Random Access Memory (DRAM) where the silicon for the stacked transistors is produced using selective epitaxy technology or laser recrystallization. This concept may be unsatisfactory as the silicon processed in this manner may have a higher defect density when compared to single crystal silicon and hence may suffer in performance, stability, and control. It may also require higher temperatures than the underlying metallization or low-k intermetal dielectric could be exposed to without reliability concerns.\nSang-Yun Lee in U.S. Pat. No. 7,052,941 discloses methods to construct vertical transistors by preprocessing a single crystal silicon wafer with doping layers activated at high temperature, layer transferring the wafer to another wafer with preprocessed circuitry and metallization, and then forming vertical transistors from those doping layers with low temperature processing, such as etching silicon. This may be less than satisfactory as the semiconductor devices in the market today utilize horizontal or horizontally oriented transistors and it would be very difficult to convince the industry to move away from the horizontal. Additionally, the transistor performance may be less than satisfactory as a result from large parasitic capacitances and resistances in the vertical structures, and the lack of self-alignment of the transistor gate.\nA key technology for 3D IC construction may be layer transfer, whereby a thin layer of a silicon wafer, called the donor wafer, may be transferred to another wafer, called the acceptor wafer, or target wafer. As described by L. DiCioccio, et. al., at ICICDT 2010 pg 110, the transfer of a thin (about tens of microns to tens of nanometers) layer of mono-crystalline silicon at low temperatures (below approximately 400° C.) may be performed with low temperature direct oxide-oxide bonding, wafer thinning, and surface conditioning. This process is called “Smart Stacking” by Soitec (Crolles, France). In addition, the “SmartCut” process is a well understood technology used for fabrication of SOI wafers. The “SmartCut” process employs a hydrogen implant to enable cleaving of the donor wafer after the layer transfer. These processes with some variations and under different names may be commercially available from SiGen (Silicon Genesis Corporation, San Jose, Calif.). A room temperature wafer bonding process utilizing ion-beam preparation of the wafer surfaces in a vacuum has been recently demonstrated by Mitsubishi Heavy Industries Ltd., Tokyo, Japan. This process allows room temperature layer transfer."}
{"text": "1. Field of the Invention\nThe present invention relates to data problem diagnosing tools and, more particularly, to a method and system for systematically diagnosing data problems in a database.\n2. Description of the Related Art\nDataBase Administrators (DBAs) monitor and maintain databases to ensure accurate and consistent storage of data in the databases. DBAs load and retrieve data to and from the database using a programming language called SQL (Structured Query Language) which is well known to DBAs. In a database, data are stored in tables. Each table has a set of rows and fields (columns). Each row in the table is identifiable by one or a combination of unique identifiers known as “keys” and each field in the table is identifiable by a field name. A field containing the keys is known as a “key field.” DBAs and database programmers utilize keys and field names to retrieve particular data from a database table.\nConventionally, when the integrity of data stored in a database is questioned, the DBA attempts to diagnose the nature of the problem and/or the exact location of the problem using one of the standard SQL statements known as “SELECT.” A SELECT statement is an SQL command for retrieving data from one or more tables of the database. Without knowing the exact nature or location of the data problem in the database, the DBA must speculate on where the problem may lie and use the SELECT statement to retrieve data from the potential problem areas of the database. The DBA compares the retrieved data with some source to determine whether inaccurate or inconsistent data are stored in the database, and whether any data is missing from a particular location in the database. This process is repeated until the DBA can diagnose properly the data problem in the database. As a result, the conventional process of diagnosing data problems in the database can be extremely time consuming and inefficient. Furthermore, the DBA must keep track of any data retrieved during this process and manually compare data sets to determine the exact nature and/or location of the data problem. This also can be tedious and is prone to human error.\nTherefore, there is a need for an improved technique by which data integrity problems or other problems in a database can be diagnosed more quickly and more systematically, thereby overcoming problems encountered in conventional data problem diagnosing techniques for databases."}
{"text": "1. Field of the Invention\nThe present invention relates generally to microelectronic structures and fabrication methods, and more particularly to the formation of interconnect insulation having low dielectric constants.\n2. Background\nAdvances in semiconductor manufacturing technology have led to the development of integrated circuits having multiple levels of interconnect. In such an integrated circuit, patterned conductive material on one interconnect level is electrically insulated from patterned conductive material on another interconnect level by films of material such as silicon dioxide.\nA consequence of having of patterned conductive material separated by an insulating material, whether the conductive material is on a single level or multiple levels, is the formation of undesired capacitors. The parasitic capacitance between patterned conductive material, or more simply, interconnects, separated by insulating material on microelectronic devices contributes to effects such as RC delay, power dissipation, and capacitively coupled signals, also known as cross-talk.\nFIGS. 1(a)-(b) illustrate a conventional approach to providing an insulating material between interconnect lines. A layer of conductive material is deposited onto the surface of a substrate 102 and then patterned to provide interconnect lines 104, as shown in FIG. 1(b). An inter-layer dielectric 106 is then formed as shown in FIG. 1(b). Inter-layer dielectric 106 generally fills the space between and above interconnect lines 104. The parasitic capacitance seen by an interconnect line is a function of the distance to another conductor and the dielectric constant of the material therebetween.\nOne way to reduce the unwanted capacitance between the interconnects is to increase the distance between them. Increased spacing between interconnect has adverse consequences such as increased area requirements, and corresponding increases in manufacturing costs. Another way to reduce the unwanted capacitance between the interconnects is to use an insulating material with a lower dielectric constant.\nWhat is needed is a structure providing low parasitic capacitance between patterned conductors, and methods of making such a structure."}
{"text": "In the field of broadcasting system, a satellite broadcasting system which can provide with local service covering a wide range has been developed and diversification of broadcasting service is being planned. In such a satellite broadcasting system, compared with an existing broadcasting system, without requiring any large-scale broadcasting equipment on the ground, it is realizable that broadcasting signals containing images and sounds are provided to a wide range of service area.\nHowever, conventional satellite broadcasting system has a problem that it cannot implement reception unless it comprises comparatively large-scale reception system. That is, radio wave for satellite broadcasting is very weak, and therefore was considered to require a comparatively large receiving antenna, etc.\nRadio wave from a satellite is also being received by a base station, etc. so that this is redistribution via a CATV network, but a reception terminal (set top box) of CATV, etc. is necessary. In addition, under a mobile environment it does not work.\nThis invention was made in view of the above described actual situation, and a first objective hereof is to make digital contents transmittable to a mobile body with a small scale deployment.\nAs a conventional contents distribution apparatus for such a mobile body terminal, for example, “Information Supply system for Vehicles” (Japanese Laid Open Gazette No. 10-73440) which is configured to receive music or map data by connecting a plug for information transmission provided in the vicinity of an oil supply nozzle at a gas station with a connector for reception provided in the vicinity of an oil supply orifice of an automobile, “Digital Information Distribution System” (Japanese Laid Open Gazette No. 11-168464) which receives distribution of music information and the like by connecting an information recording-reproducing device in a customer's possession with a vending machine exclusive for digital information, “Digital Information Distribution System” (Japanese Laid Open Gazette No. 11-191266) which is configured so that a plurality of music distribution terminals are brought into connection with a service center which digitalize a large amount of music data to store by way of communication lines, and “Digital Information Distribution System” (Japanese Laid Open Gazette No. 2000-48479) which is configured to perform reproduction by bringing a user's information reproducing terminal into connection with a music information vending terminal via a predetermined interface, sending and receiving encoded music information and the user's ID, and decoding the encoded information, and the like, are known.\nNext, vehicle electronic equipment as a receiving terminal installed in a mobile body in the above described contents distribution system is provided with various antitheft device.\nFor example, panel detachable type vehicle electronic equipment (panel detachable type equipment), in which an operation panel disposed in the front surface is detachable, is given as an example. This panel detachable equipment has a operation panel which is removed by a user at the time when he/she leaves the car for a long time so as to make the equipment unusable as well as to place the equipment under incomplete state for antitheft. An example of technology of such a panel detachable equipment is disclosed in Japanese Laid Open Gazette No. 6-252565.\nIn such conventional panel detachable equipment, antitheft can be planned by physically removing a operation panel, but the operation panel under the removed state could not be used effectively. That is, the operation panel under the removed state could not be used at all, and in spite of its smallness, it turned out to be burdensome.\nA similar problem is broadly applicable to electric equipment comprising a detachable operation panel.\nThe present invention, which was made in view of the above described actual conditions, has a second objective to provide electric equipment which can utilize effectively a operation panel under removed state, and a panel which can be utilized effectively under removed state, and to provide vehicle electronic equipment for planning antitheft as well as effective use.\nMoreover, as the above described electric equipment, there is one which uses a semiconductor element, especially a flash memory which is a nonvolatile memory having a large capacity and being rewritable, as storage medium. This kind of electric equipment does not need any rotary driving portion, and can make an apparatus smaller and control consumed energy.\nHowever, a flash memory has a problem that it has a narrow working temperature range. Normally, around 65° C. is the limit. In addition, when the temperature to not less than the temperature, recorded data could even disappear.\nIn normal electric equipment, the temperature inside the apparatus does not frequently up to 65° C., and even so, can be coped with an appropriate ventilating mechanism being provided.\nHowever, in case of an apparatus to be mounted on a vehicle, due to its peculiarity on installation site such as that the main body portion is stored in a dashboard of the vehicle, it is not exceptional that the temperature increases to not less than 65° C. Moreover, inside a vehicle in summer, even at time of stopping, the temperature happens to reach around 100° C.\nDue to this, simple use of a flash memory in electric equipment to be mounded on a vehicle as memory media gives rise to a problem that it lacks stability and reliability.\nSimilar problem commonly exists in electric equipment to be mounded on vehicles which stores not only music data but also image data such as map data, text data, etc. in a flash memory for use.\nThe present invention, which was made in view of the above described actual conditions, has a third objective to improve reliability of electric equipment to be mounded on a vehicle in which semiconductor recording element is used as storage medium, and to use the semiconductor recording element within an appropriate temperature range."}
{"text": "The task of such linking, looping or hemming machines is in particular to cleanly and durably link or loop the edges of carpets or the like in the form of a roll product and which are cut to the use size. The aim of the invention is that despite the relatively high capacity of the machine, it can be easily handled and can also be used in a mobile manner, because in particular when looping very large surfaces the space required for moving the complete material portion passed the machine would be excessive.\nThe machine according to the invention has a drive, a sewing mechanism with means for conveying the material relative to the sewing mechanism, a needle mechanism and a loop forming mechanism, as well as a thread feed to the needle and loop forming mechanisms.\nAccording to one feature of the invention the sewing mechanism preferably constructed for forming a two-thread overlock or overcast seam has a retaining gripper, which holds back a thread loop during the formation of an overlock or overcast loop. It preferably engages close to the needle movement path from below into the thread coming out of a yarn gripper.\nU.S. Pat. No. 4,062,307 discloses a sewing machine producing a single-thread overcast seam. It uses a yarn gripper gripping in a loop the thread passed through the material with the needle and is passed around the edge for the next needle perforation in said loop. This machine operates in a completely satisfactory manner and can also loop thick carpets with a limited constructional size and high capacity. However, relatively large thread lengths must be drawn at high speed through the needle hole, because the entire thread supply or feed takes place from above via the needle. As the thread for looping normally consists of a relatively slightly turned thick and bulky yarn, in the case of floor coverings with a very rigid and firm back working with great care is necessary, so that the yarn does not unravel and tear during the drawing through the hole."}
{"text": "With the recent digitization of communications, performance improvement such as enhancement in bit resolution and increase in conversion speed has been increasingly demanded of A/D converters (ADCs) used in the field of digital communications. With improvement in the performance of ADCs, however, the power consumption thereof often increases. For example, when it is attempted to increase the conversion speed of a sampling ADC, a large current must be fed to permit high-speed charge/discharge of capacitance elements for sampling of an input signal. Considering application of ADCs to mobile equipment such as cellular phones, power reduction will also be necessary in ADCs as well as the performance improvement.\nAs a high-performance, low-power ADC, a column parallel ADC (column ADC) composed of super-many sampling ADCs is known (see Non-Patent Document 1, for example). In a column ADC, in which hundreds to thousands of sampling ADCs operate in parallel, very high-speed A/D conversion capability can be attained as a whole even though a single sampling ADC is slow in operation.    Non-Patent Document 1: Yoshikazu Nitta et al., “High-Speed Digital Double Sampling with Analog CDS on Column Parallel ADC Architecture for Low-Noise Active Pixel Sensor”, ISSCC 2006/SESSION 27/IMAGE SENSORS/27.5"}
{"text": "Printed circuit boards (PCBs) denote boards just before electric components are mounted. In such a PCB, a circuit line pattern is printed on an insulation board using a conductive material such as copper. That is, the PCB denotes a circuit board in which an installation position of each of the components is decided and a circuit pattern connecting the components to each other is printed and fixed on a surface of the flat plate to densely mount various electric devices on the flat panel.\nSpecifically, a multi-layer PCB may be used to be built in a cellular phone, a video camera, a notebook, etc., in which high-integration and compact size are required in recent years. The multi-layer PCB (MLB) is manufactured by building up the PCB one by one. In the build-up process, the board may be manufactured and evaluated one by one to improve yield of the multi-layer PCB. Also, the layers may be precisely connected to each other using a wire to realize a small-sized PCB. Hereinafter, a PCB buried within the multi-layer PCB will be referred to as an inner PCB.\nFor example, in case of a PCB including ten layers, eight inner PCBs may be provided.\nAlso, an insulation layer may be disposed between the multi-layer PCBs. A circuit pattern metal part, a via hole metal part electrically connected to the inner PCB, and a pad metal part disposed on the via hole metal part may be disposed on a surface of the insulation layer. A trench process using a laser may be performed as a process of forming engraved patterns (i.e., circuit pattern, via hole, and pad part) in which the metal part, the via hole metal part, and the pad metal part are disposed in the insulation layer."}
{"text": "This invention relates to a method and apparatus for recording audio information in an authenticatable, tamper-proof manner.\nTraditionally, written documents have been used to provide permanent records of transactions and agreements. One example of this type of document is a contract for the sale of an item, which typically identifies the name of the parties, the date, the subject matter of the contract, and a price. The contract provides a permanent record that can be used at a later date to establish the terms of the agreement between the parties.\nOral contracts, on the other hand, do not provide a permanent record of the terms of the agreement. As a result, if a dispute arises over the terms of the agreement at a later date, it becomes difficult to prove exactly what the parties agreed to, or whether they made a binding contract at all. Because there is no permanent record, an unscrupulous party could be untruthful about the agreed-upon terms to escape his obligations. Even absent dishonesty, parties to an oral contract may have different recollections of exactly what they agreed to. Moreover, one of the persons who entered into the agreement may be permanently or temporarily unavailable. These problems tend to worsen as time passes.\nBecause of these problems, all states have statutes declaring that certain oral agreements are unenforceable, typically including the sale of land, and the sale of goods exceeding a certain value. If a trustworthy record of an oral agreement or transaction could be obtained, however, the problems of oral agreements could be overcome.\nExisting methods of recording conversations, however, do not address these problems. For example, telephone answering machines, tape recorders, and handheld digital audio recording devices can be used to record a voice or a conversation. It is, however, relatively easy to delete or to alter the recorded audio information. In particular, readily available electronic devices can splice sections out of an audio conversation, and can even rearrange words to make it appear that a party said something that was never actually said. Moreover, there is no effective way for parties to sign an audio recording. As a result, it may be difficult to identify the parties that actually agreed to the terms contained in an audio conversation and intended to abide by such terms. Further, there is no way known to applicants to verify that an oral negotiation matured into an agreement.\nIn addition, existing telephone answering machines and tape recorders do not provide a reliable indication of when the conversation occurred. While some answering machines do record the time a call was received, this “time stamp” is extremely unreliable because a party could rerecord a new time over the time recorded by the answering machine. Alternatively, a party might either intentionally or accidentally set the date on an answering machine incorrectly. This would allow two corroborating parties to pretend that they made an agreement on a certain date, even though the agreement was not made until a later date. As a result, telephone answering machines and ordinary cassette recorders do not alleviate the problems of oral agreements described above.\nSTEN-TEL is an example of a system designed specifically for recording telephonic audio information. STEN-TEL is available from Sten-tel Inc. (having a place of business at 66 Long Wharf, Boston, Mass. 02110). To use STEN-TEL, a person places a telephone call to the STEN-TEL server, and the server digitally records the telephone call. After the digital recording is made, a transcriptionist accesses the recording and generates a typed record of the telephone call. The typed transcription is then uploaded to the server, where it is stored. Permanent storage of the digitally recorded audio conversation is optional. After the transcription is stored in the server, it can be downloaded to the users. Every transcription is assigned a unique identification number, and all status information is maintained in a centralized database.\nThe STEN-TEL system does not, however, overcome the drawbacks of existing telephone answering machines and audiocassette recorders. First, the ability to restrict access to files is limited or non-existent. Apparently any person who has the file identifier can access the stored information. Second, the information is vulnerable to tampering. Third, although STEN-TEL apparently stores the time of the call, time stamps are not embedded into the stored information. This makes STEN-TEL vulnerable to modifications of the stored date for a given conversation. Finally, digital signatures are not used to provide security and/or authenticate the parties.\nOne system that does incorporate certain security features is described in U.S. Pat. No. 5,594,798 (Cox et al.), which describes a voice messaging system. In Cox's system, however, an encryption key is stored along with the encrypted message. Because a hacker could obtain access to the encrypted message by retrieving the encryption key, Cox's system is vulnerable to attack. In addition, Cox's system is intended for use with secure telephone devices (STD). Ordinary telephones cannot call into Cox's system to have an audio message recorded.\nNo existing audio recording system is known to applicants that facilitates the permanent recording of an audio conversation in an authenticatable form so that a user can simply place a telephone call to a central server and have the server encrypt the conversation and record the time of the conversation, all in a tamper-proof manner."}
{"text": "The present invention relates to ultrahigh molecular weight polyethylene and polypropylene fibers having high tenacity, modulus and toughness values and a process for their production which includes a gel intermediate.\nThe preparation of high strength, high modulus polyethylene fibers by growth from dilute solution has been described by U.S. Pat. No. 4,137,394 to Meihuizen et al. (1979) and pending application Ser. No. 225,288 filed Jan. 15, 1981.\nAlternative methods to the preparation of high strength fibers have been described in various recent publications of P. Smith, A. J. Pennings and their coworkers. German Off. No. 3004699 to Smith et al. (Aug. 21, 1980) describes a process in which polyethylene is first dissolved in a volatile solvent, the solution is spun and cooled to form a gel filament, and finally the gel filament is simultaneously stretched and dried to form the desired fiber.\nUK patent application GB No. 2,051,667 to P. Smith and P. J. Lemstra (Jan. 21, 1981) discloses a process in which a solution of the polymer is spun and the filaments are drawn at a stretch ratio which is related to the polymer molecular weight, at a drawing temperature such that at the draw ratio used the modulus of the filaments is at least 20 GPa. The application notes that to obtain the high modulus values required, drawing must be performed below the melting point of the polyethylene. The drawing temperature is in general at most 135.degree. C.\nKalb and Pennings in Polymer Bulletin, vol. 1, pp. 879-80 (1979) J. Mat. Sci., vol. 15, 2584-90 (1980) and Smook et al. in Polymer Bull., vol. 2, pp. 775-83 (1980) describe a process in which the polyethylene is dissolved in a nonvolatile solvent (paraffin oil) and the solution is cooled to room temperature to form a gel. The gel is cut into pieces, fed to an extruder and spun into a gel filament. The gel filament is extracted with hexane to remove the paraffin oil, vacuum dried and then stretched to form the desired fiber.\nIn the process described by Smook et al. and Kalb and Pennings, the filaments were non-uniform, were of high porosity and could not be stretched continuously to prepare fibers of indefinite length."}
{"text": "A data cube, such as a Microsoft® SQL Server Analysis Services cube, is a three-dimensional (or higher) database structure for storing and presenting data. Users of such a cube want their data to be highly available, with minimum interruptions in the event of hardware or software faults. However, there was heretofore not any consistent way to ensure highly available cubes, including highly available cubes in which data access remains efficient, even in heavy traffic."}
{"text": "Technical literature has disclosed mortars for either interior or exterior plaster possessed of either anticondensate or animoisture diffusive, heat-insulating, or biocidal properties.\nThe main disadvantages of these mortars is that they do not concentrate all the above-mentioned properties in one single product, that they consist of expensive materials (silicones, palmitates, etc.), and that they necessitate laborious preparation and application techniques in use."}
{"text": "1. Field of the Invention\nThe present invention relates to a connecting structure of a fender apron, and more particularly, to a connecting structure for a fender apron that can disperse an impact load when a vehicle collides.\n2. Description of Related Art\nWhen a vehicle collides head-on, a front side member of the vehicle serves to absorb an impact.\nHowever, according to an investigation, accidents in which the vehicle collides off the front side member occupy 25% of actual accidents in a collision accident.\nIn the accidents in which the vehicle collides off the front side member, a fender apron serves to transfer an impact load.\nA general collision load transferring component is constituted by three units of an under body unit, for example, a front apron member, a side body unit, for example, an A-pillar upper reinforcement unit and a door hinge reinforcement unit, and a cowl unit, for example, a cowl side member.\nHowever, in the case of coupling of the respective units, respective connection sections are not continuous, and thus coupling force is weak, and the respective units are connected to each other by a flange, and as a result, a process for connection needs to be added and a deformation possibility by impact load transferring is high.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
{"text": "A snowboard is a long continuous surface platform made from a variety of materials designed to capture certain aspects of surfing on snow. When using a snowboard, there is a sensation of gliding over a surface and shifting one's body weight from one side of the board to another in order to execute a turn in either direction. In surfing, the execution and completion of a turn relies on the surface hull design (single concave, double concave), fins (single, double, triple), and overall body design (teardrop, asymmetrical). The execution of a turn with a snowboard is based on the flex and shape of the bow of the board and the ratio of the width of the bow to the width of the waist. The completion of the turn is based on a mixture of the flex and shape of the tail and the ratio of the width of the tail to the width of the waist. Ideally, a turn with a snowboard is initiated by applying pressure to the downhill foot and leaning down and into the side of the board one wishes to turn.\nSkateboarding utilizes an articulating platform with wheels attached to trucks mounted to the underside of a platform that invariably is designed to appear as a surfboard. The execution and completion of a turn with a skateboard is accomplished by shifting weight from one side of the platform to the other and maintaining pressure slightly to the bow. The articulating wheels (front set turns one direction while the back set turns another) allow for completion of the turn. With a skateboard, the turns can be of varying radius and frequency.\nSnowboarding and skateboarding attempt to capture aspects of surfing. Because of the snowboard's inherent design limitations it does not attain certain performance parameters in the hands of the average user. Specifically, short radius turns and high-frequency edge-to-edge turns are difficult and not attainable for the recreational user. In addition, the time to learn how to use a snowboard can be long and frustrating, causing some users to avoid attempts to learn. Control of the board is essential and time consuming to master. As noted above, control relies on shifting weight and movement of the uphill foot, from side to side, to assist in the turn cycle. It would be desirable to provide a system for guiding an apparatus on a surface, such as snow, that provides the sensation of surfing, i.e., leaning from side to side to carry out a turn, as well as edge-to-edge control that allows a user to achieve a sensation of cutting up and down the face of a wave while minimizing the loss of vertical feet. Such a device would desirably allow the first-time user to readily master the necessary skills which would further promote usage of the device."}
{"text": "This invention relates to microwave power combiners using solid state amplifier devices and to planar transitions from waveguide to microstrip or vice versa.\nCombining the power from several FET amplifiers is a current approach for obtaining higher power microwave solid state power sources. The combining efficiency must be increased and the fabrication costs must be reduced to make this approach competitive with conventional vacuum tube power sources such as traveling-wave tubes. The successful solution will simultaneously satisfy all the important requirements of low cost, high combining efficiency and broad bandwidth, and does not use coaxial connectors.\nA fin-line transition from waveguide to microstrip is inherently low loss. A single transition in a single waveguide is described in J. H. C. van Heuven, IEEE Transactions on Microwave Theory and Techniques, Vol. MTT-24, No. 3, March 1976, pp. 144-147."}
{"text": "With the advanced computing technologies available today, more and more people demand instant access to data and other information. Further, since such a vast amount of data is readily available, people tend to become more inquisitive and search for information on a large variety of topics and/or locations. At times, the desire and/or need for information can be overwhelming and, due to time constraints, commitments, and other limitations (e.g., putting a task off and forgetting about it), people might not have time to perform all the research or obtain all the desired information.\nMany people utilize avatars as a way to express a computer generated representation of themselves or as their alter ego. The avatar can be represented as a two-dimensional icon (e.g., picture) or a three-dimensional model. Generally, avatars are constructed to represent a friend or assistant who can interact with the user or an environment of the user as a distinct entity in relation to the perspective of that user. However, today's avatars are arbitrarily discussed or implemented as distinct third-person entities. Once implemented these avatars generally do nothing more than visually represent an entity."}
{"text": "Non-volatile data storage devices, such as embedded memory devices (e.g., embedded MultiMedia Card (eMMC) devices) and removable memory devices (e.g., removable universal serial bus (USB) flash memory devices and other removable storage cards), have allowed for increased portability of data and software applications. Users of non-volatile data storage devices increasingly rely on the non-volatile storage devices to store and provide rapid access to a large amount of data.\nSome data storage devices may include volatile memory, and data stored in the volatile memory may be lost in the case of power loss. In some data storage devices, a single-bit per cell (also known as single-level cell (SLC)) memory may be used to store data that can persist in the case of power loss. However, SLC memory may have high cost and limited endurance (e.g., data retention)."}
{"text": "1. Field of the Invention\nThis invention generally relates to an apparatus for recording digital information on a recording medium such as a magnetic disk or tape and reproducing the recorded information and more particularly to a code-error correcting device for correcting code-errors in digital signals used in the apparatus for recording digital information on the recording medium and reproducing the recorded information.\n2. Description of the Related Art\nReferring first to FIG. 8, there is shown the structure of a code of digital signals used in a typical apparatus for recording audio signals on a magnetic tape by means of a rotary head and reproducing the recorded information (that is, what is called an R-DAT (Digital Audio Tape recorder)) therefrom. As shown in this figure, the code includes data (DATA) composed of 28.times.26 symbols, a transverse or horizontal parity code (C.sub.2 PARITY) composed of 28.times.6 symbols and a longitudinal or vertical parity code (C.sub.1 PARITY) composed of 4.times.32 symbols. In the case of Reed Solomon Code (R.S.C), sets of data concerning the parity codes C.sub.1 and C.sub.2 are (32, 28, 5) and (32, 26, 7), respectively. In each of the parentheses, a first, second and third numeral indicates values of the total length of a code, the length of data and a minimum distance between code words, respectively.\nFurther, referring now to FIG. 9, there is shown the format of signals employed when recording the signals having such structure of codes. In this figure, reference characters SYNC indicates a synchronizing signal; ID an identification signal; ADR an address signal; P a block parity signal; DATA data of 28 symbols; and C.sub.1 a C.sub.1 parity code of 4 symbols. That is, signals SYNC, ID, ADR and P are added to data signals. Incidentally, in this case, the block parity signal is given by EQU P=ID.sym.ADR.\nNamely, the signals, of which the format is as shown in FIG. 9, are recorded on the magnetic tape and reproduced therefrom.\nThe above described conventional apparatus can detect address errors to some extent by transmitting the block parity signal indicating ID.sym.ADR together with the signal indicating data DATA. However, the conventional apparatus has a drawback that the capability of detecting the address errors is not sufficient to precisely detect the address error and as a consequence the address errors increase. In this case, data are stored in an erroneous area within a memory in accordance with the erroneous address information because the address information generally determines an area in the memory in which data are to be stored. Conventionally, even when the error cannot be detected by using the longitudinal parity code (C.sub.1), the error can still be corrected if the error is present within the range which can be corrected by using the transverse parity code (C.sub.2). Further, if the error exceeds the capability of detecting the error by using the transverse parity code (C.sub.2), it is necessary to locate the error on the basis of the error information which is generated after the check by using the parity codes C.sub.1 and C.sub.2.\nIn such case, if only the area in the memory is erroneous and the parity code C.sub.1 is correct, there is the inconvenience that in spite of the fact that a sequence of data is erroneous, the error cannot be detected. Thus, to eliminate such inconvenience, it has been proposed that when the parity code C.sub.1 is generated, the address is included as a generating element for the parity code C.sub.1. Such approach has a defect that the capability of correcting error is degraded because the code is not a product code.\nTherefore, it is an object of the present invention to provide a code correcting device of which the capability of correcting error is significantly improved, thereby decreasing the possibility of passing over the error."}
{"text": "In a computerized content delivery network, content providers typically design and provide content items for delivery to a user device via one or more content slots of an electronic resource. Content items typically include display content items and textual content items. Display content items often include images, video, graphics, text, or other visual imagery. Textual content items often include a headline and a short “creative” text. It can be challenging for content providers to create effective and attractive display content items. It can also be challenging for content providers to write the creative text portion of textual content items.\nVarious templates and stock elements have been used to partially automate the process of creating content items. However, content items created from rigid templates and stock elements are often stale, unattractive, and not well-suited to the particular business, product, or service featured in the content item."}
{"text": "1. Field of the Invention\nThis invention relates to disk drive suspensions, and, more particularly, to a disk drive suspension comprising a distal portion, a base portion and a spring portion between the distal portion and the base portion, and a stiffener attached to the suspension to lend a desired stiffness to the suspension. The suspension spring portion typically exerts a predetermined gram load on the distal portion to maintain a slider carried by the distal portion into a desired operating proximity to a rotating disk. The invention specifically relates to the addition of a second spring complementary in function to the suspension spring. The second spring traverses the suspension spring portion and may connect separated sections of the stiffener. The second spring adds a highly controllable spring force to the suspension to better control of the suspension spring rate, e.g. when the suspension material thinness limits spring rate obtainable from just the suspension, to reduce gram losses with use, and for other purposes to be described.\n2. Description of the Related Art\nSuspensions typically comprise a metal member that can be a load beam having a base portion, a spring portion and a rigid distal portion over which a conductor is passed to connect a slider carried by a separate flexure on the load beam distal portion to the device electronics, or the metal member can be the metal layer in a laminate of the metal layer, an insulative layer, and trace conductors as seen in so-called wireless suspensions. Wireless suspension metal layers are very thin and are frequently provided with supplemental stiffening. Stiffeners are added to the base portion and/or to the distal portion of the metal layer."}
{"text": "Cotton is an important fiber crop in many areas of the world. The methods of biotechnology have been applied to cotton for improvement of the agronomic traits and the quality of the product. The method of introducing transgenes into cotton plants has been demonstrated in U.S. Pat. No. 5,004,863. One such agronomic trait important in cotton production is herbicide tolerance, in particular, tolerance to glyphosate herbicide. This trait has been introduced into cotton plants and is a successful product now used in cotton production. The current commercial Roundup Ready® cotton event (1445) provides excellent tolerance to glyphosate, the active ingredient in Roundup®, through the four-leaf stage (Nida et al., J. Agric. Food Chem. 44:1960-1966, 1996; Nida et al., J. Agric. Food Chem. 44:1967-1974, 1996). However, foliar application beyond the four-leaf stage must be limited due to insufficient tolerance in male reproductive tissues in certain environmental conditions. This lack of male reproductive tolerance appears to be a result of insufficient CP4 EPSPS expression in critical tissues, higher sensitivity of these tissues to glyphosate, and accumulation of high amounts of glyphosate in these strong sink tissues (Pline et al., Weed Sci. 50:438-447, 2002). There is a need for a cotton plant more highly glyphosate tolerant than Roundup Ready® cotton 1445.\nIt would be advantageous to be able to detect the presence of a particular event in order to determine whether the progeny of a sexual cross contain a transgene of interest. In addition, a method for detecting a particular event would be helpful for complying with regulations requiring pre-market approval or labeling of foods derived from recombinant crop plants, for example. It is possible to detect the presence of a transgene by any well known nucleic acid detection method such as the polymerase chain reaction (PCR) or DNA hybridization using nucleic acid probes. These detection methods generally focus on frequently used genetic elements, such as promoters, 3′ transcription terminators, marker genes, etc. As a result, such methods may not be useful for discriminating between different events, particularly those produced using the same DNA construct unless the sequence of genomic chromosomal DNA adjacent to the inserted DNA (“flanking genomic DNA”) is known. Event-specific DNA detection methods for a glyphosate tolerant cotton event 1445 have been described (US 20020120964, herein incorporated by reference in its entirety).\nThe present invention relates to a glyphosate tolerant cotton event MON 88913, compositions contained therein, and to the method for the detection of the transgene/genomic insertion region in cotton event MON 88913 and progeny thereof."}
{"text": "(1) Field of the Invention\nThe disclosure herein relates to a solar cell and a method of fabricating the solar cell.\n(2) Description of the Related Art\nSolar cells use a photovoltaic effect to convert optical energy into electrical energy. Solar cells may be classified into silicon solar cells, thin film solar cells, dye-sensitized solar cells, and organic polymer solar cells, according to their materials. Solar cells are used as main or sub power sources for various electronic products, artificial satellites, and rockets.\nSuch a solar cell includes: a first semiconductor having a first conductive type, a substrate; and a second semiconductor disposed on the front surface of the first semiconductor and having a second conductive type opposite to the first conductive type. The first and second semiconductors form a p-n junction. Light is shed on the p-n junction to form pairs of electrons and holes in the first and second semiconductors. The holes and the electrons are moved by an electrical potential difference in the p-n junction, thereby generating electric current."}
{"text": "I. Field of the Invention\nThe present invention relates to a connector plate for securing abutting wooden truss members together.\nII. Description of the Prior Art\nThe use of nail connector plates for securing together abutting truss members is well known in the art. These previously known connector plates typically comprise a sheet of metal having a plurality of outwardly extending prongs formed on one side of the sheet which are hammered into the abutting truss members.\nIt has been the conventional practice with many of these previously known connector plates to nail the connector plate to the wooden truss elements after the insertion of the prongs into the truss elements. The nailed attachment of the connector plate to the truss elements prevents unintended detachment of the connector plate from the truss members. The nail attachment is disadvantageous, however, in that it requires additional materials, i.e. the nails, and is more time consuming to construct, and therefore, more expensive in labor costs.\nDue to the disadvantages of nailing the connector plate to the truss element, there have been other previously known connector plates which are secured to the truss element only by hammering the prongs into the truss elements. These previously known connector plates, such as described in U.S. Pat. No. 3,603,197, issued Sept. 7, 1971, typically include a barb on each point prong which bites into the truss elements after insertion of the prongs. Connector plates of this type, however, are disadvantageous for a number of reasons.\nFirst, due primarily to the sharp tip on each prong, after long usage the prongs tend to slide outwardly from their hole and away from the truss elements. This partial disengagement of the prongs from their respective holes greatly weakens the truss element joint.\nThe previously known connector plates of this type are also disadvantageous in that the connector plates are difficult to handle due to the sharply pointed end on the prongs. Injuries to workmen handling the connector plates is common in the trade, even when the workmen are wearing gloves.\nLastly, these previously known connector plates are disadvantageous in that a complex punch and die arrangement is required to produce the barbs on the prongs during the punching operation. Due to their resultant high cost, the required punch and die arrangements for these previously known connector plates unduly increases the overall cost for the individual connector plates."}
{"text": "In the field of electromagnetic sensing, certain sensing structures place strict requirements on the polarization of the electromagnetic radiation for absorption by the structure. To meet these requirements, a diffraction grating is used to convert the electromagnetic radiation into a polarization or mode which is absorbed by the sensing structure. The inclusion of a reflector in the sensor creates an electromagnetic cavity. When the incident radiation is at a frequency which resonates within the cavity, a standing wave results, creating regions of high and low electric field intensities. The magnitude of the detection signal is proportional to the number and magnitude of the high intensity field regions. The boundary conditions within the electromagnetic cavity have a significant impact on the magnitude of the detection signal. Improper boundary conditions cause the electromagnetic field to interfere destructively with itself, diminishing the number and magnitude of high field regions, especially near the periphery of the sensing structure.\nOne such electromagnetic sensor is the quantum well infrared photodetector (QWIP). With current objectives to develop large area focal plane array (FPA) technology for mid wave-length infrared radiation (MWIR), long wavelength infrared radiation (LWIR), and multi-spectral applications at low cost with high performance, QWIP technology is being extensively explored. QWIP technology suffers some performance disadvantage relative to other infrared (IR) technologies. In view of the more mature material and processing technology utilized with QWIP FPAs relative to other IR technologies, there exists a need for design improvements to enhance quantum efficiency (QE) and therefore performance."}
{"text": "This invention relates generally to refrigeration and operation and more particularly to a method and apparatus for boosting the cooling capacity and efficiency of air-conditioning systems under a wide range of ambient atmospheric conditions.\nIn air conditioning, the basic circuit is essentially the same as in refrigeration. It comprises an evaporator, a condenser, an expansion valve, and a compressor. This, however, is where the similarity ends. The evaporator and condenser of an air conditioner will generally have less surface area. The temperature difference .DELTA.T between condensing temperature and ambient temperature is usually 27.degree. F. with a 105.degree. F. minimum condensing temperature, while in refrigeration the difference .DELTA.T can be from 8.degree. F. to 15.degree. F. with an 86.degree. F. minimum condensing temperature.\nI have previously improved the cooling capacity and efficiency of refrigeration systems. As disclosed in my U.S. Pat. No. 4,599,873, this is accomplished by addition of a liquid pump at the outlet of the receiver or condenser. Operation of the pump adds 5-12 p.s.i. of pressure to the condensed refrigerant flowing into the expansion valve, a process I call liquid pressure amplification. This suppresses flash gas and assures a uniform flow of liquid refrigerant to the expansion valve, substantially increasing cooling capacity and efficiency. The best results are obtained when such a system is operated with the condenser at moderate ambient temperatures, usually under 80.degree. F. As ambient temperatures rise above the minimum condensing temperature, the advantages gradually decrease. The same thing happens when the principles of my prior invention are applied to air conditioning, except that the minimum condensing temperature is higher.\nWhile conventional air-conditioning systems can benefit from my prior invention, the greatest need for air conditioning is when ambient temperatures are high, over 80.degree. F. Conventional air conditioning becomes less effective and efficient as ambient temperatures rise to 100.degree. F. or more, as does use of my prior liquid refrigerant pressure amplification technique.\nIt is, therefore, an object of the invention to improve the efficiency of refrigeration and air-conditioning systems.\nAnother object of the invention is to increase the cooling capacity of such systems when operated at high ambient temperatures.\nA further object of the invention is to enable the aforementioned objects to be attained economically and by retrofitting existing systems as well as in new systems.\nThe present invention is an improvement in the structure and method of operation of an air-conditioning or refrigeration system which includes a compressor, a condenser, an expansion valve, an evaporator, and conduit means interconnecting the compressor, condenser, expansion valve and evaporator in series in a closed loop for circulating refrigerant therethrough, and optionally may include a receiver between the condenser and expansion valve. The conduit means includes first conduit means coupling an outlet of the compressor to an inlet to the condenser to convey superheated vapor refrigerant from the compressor into the condenser at a first pressure and temperature. A centrifugal pump means has an inlet coupled to an outlet of the condenser (or to the receiver outlet) for receiving condensed liquid refrigerant at a second pressure less than said first pressure and boosting the second pressure of the condensed liquid refrigerant by a substantially constant increment of pressure within a predetermined range to discharge the condensed liquid refrigerant from an outlet of the pump means at a third pressure greater than said second pressure. A second conduit means couples the outlet of the pump means to an inlet to the expansion valve to transmit a first portion of the condensed liquid refrigerant from outlet of the pump means at said third pressure through the expansion valve into the evaporator to vaporize and effect cooling for air conditioning or refrigeration. A third conduit means couples the outlet of the pump means to an inlet to the condenser to transmit a second portion of the condensed liquid refrigerant from outlet of the pump means into the inlet of the condenser to vaporize therein. The portion of the condensed liquid refrigerant injected into the condensor inlet cools the superheated vapor refrigerant entering the condenser to a reduced temperature, thereby reducing said first pressure.\nThe first and second conduit means are preferably proportioned so that the second portion of refrigerant is sufficient to reduce the first temperature to a reduced temperature close to a saturation temperature of the refrigerant, preferably within 10.degree. F. to 15.degree. F. above saturation temperature, and so that the second portion of refrigerant is substantially less than the first portion, preferably less than about 5% of the first portion and typically in the range of 2%-3% of the first portion. Suitably, the first and second conduit means are proportioned with a cross-sectional area ratio of about 16:1. The system preferably further includes means responsive to a temperature of the evaporator for modulating the expansion valve.\nIn the improved method of operation, superheated vapor refrigerant is transmitted from the compressor to an inlet to the condenser at a first temperature and pressure. The vapor refrigerant is condensed and discharged as liquid refrigerant at a second temperature and pressure less than said first temperature and pressure. The pressure of the liquid refrigerant discharged from the condenser (or receiver) is boosted to a third pressure greater than the second pressure by a substantially constant increment of pressure. Then, in accordance with the invention, a first portion of the liquid refrigerant is transmitted at said third pressure via the expansion valve into the evaporator and a second portion thereof is transmitted into the condenser inlet so that the first temperature of the superheated vapor refrigerant is reduced toward said second temperature, thereby reducing said first pressure.\nThe first and second portions of liquid refrigerant at said third pressure are proportioned so that the first portion is substantially greater than the second portion. Preferably, the added increment of pressure is 8 to 10 p.s.i. and the second portion has a flow rate less than 5% of the flow rate of the first portion. The flow of the first portion through the expansion valve can be modulated in response to a temperature in the evaporator.\nPrior art ammonia-refrigeration systems are known in which a portion of liquid refrigerant is injected from the receiver to the condenser inlet to suppress superheat. This has not been done, however, in combination with adding an incremental pressure, for example by means of a centrifugal pump, to the pressure of the liquid refrigerant flowing into the expansion valve.\nOperation with an added incremental liquid refrigerant pressure preferably includes allowing the first pressure to float with an ambient temperature. This reduces overall system pressures, thereby increasing system efficiency at moderate ambient temperatures. The present invention desuperheats the compressed refrigerant vapor as it enters the condenser, lowering its temperature and further reducing the first pressure, even when ambient temperatures are high. The invention thus raises the temperature range over which benefits can be obtained from adding an increment of pressure to the liquid refrigerant. This further improves efficiency and enables effective operation in very high ambient temperature environments.\nThe foregoing and other objects, features and advantages of the invention will become more readily apparent from the following detailed description of a preferred embodiment of the invention which proceeds with reference to the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of information storage devices, and, more particularly, the present invention relates to an improved system and method for allocating memory.\n2. Description of the Related Art\nIn processor-controlled devices, a memory, whether dynamic or static RAMs (Random Access Memory) is provided for the storage of programming instructions and data. In order to write data to memory areas of the memory or read data from the individual memory areas, the processor generates a memory area address and passes it to a memory controller known as an SRAM or DRAM controller to those of ordinary skill in the art. The memory controller forwards this memory area address via address lines to address inputs of the memory. The memory area addresses generally represent a binary number, the smallest addressable memory area usually comprising one or more bytes.\nThe memory is formed by a plurality of insertion regions, into which in each case at least one memory module, usually an integrated memory chip, can be inserted. The insertion regions, which are typically plug-in devices, are frequently configured such that memory modules or memory chips having different storage capacities can be inserted. For example, memory modules with one or four Kbytes or Mbytes are typically used. Consequently, because of the various insertion regions, optimum matching of the storage capacity of the memory to the respective application is possible. This means that the memory modules must be inserted into the address-conforming insertion regions as a function of their storage capacity in order to ensure error-free memory addressing. Incorrect information regarding the storage capacity and the insertion regions leads to significant disturbances in the processor system and possibly to its failure."}
{"text": "Not applicable to this application.\nNot applicable to this application.\n1. Field of the Invention\nThe present invention relates generally to baseball bat training devices and more specifically it relates to a training bat system for increasing the batting skills of a baseball player.\n2. Description of the Related Art\nBatter training devices have been in use for years. A commonly utilized batter training device is comprised of a weight having a ring structure that surrounds the barrel of a baseball bat often times referred to as a xe2x80x9cdoughnut.xe2x80x9d The weight placed upon a conventional baseball bat increases the overall weight of the baseball bat and the player then swings the baseball bat repeatedly with the weight placed upon thereof.\nWhile weights for baseball bats assist the player in developing increased strength, they do not assist the player in developing increased ball engagement accuracy. Conventional baseball bat devices do not significantly increase the mental and physical focus required to engage a baseball with the bat.\nExamples of patented devices which may be related to the present invention include U.S. Pat. No. 3,116,926 to Owen et al.; U.S. Pat. No. 6,050,908 to Muhlhausen; U.S. Pat. No. 4,682,773 to Pomilia; U.S. Pat. No. 339,621 to Briden; U.S. Pat. No. 6,280,353 to Brundage; U.S. Pat. No. 5,741,193 to Nolan; and U.S. Pat. No. 5,456,461 to Sulllivan.\nWhile these devices may be suitable for the particular purpose to which they address, they are not as suitable for increasing the batting skills of a baseball player. Conventional baseball bat training devices do not significantly assist with developing mental and physical focus for engaging a baseball.\nIn these respects, the training bat system according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in so doing provides an apparatus primarily developed for the purpose of increasing the batting skills of a baseball player.\nIn view of the foregoing disadvantages inherent in the known types of baseball bat training devices now present in the prior art, the present invention provides a new training bat system construction wherein the same can be utilized for increasing the batting skills of a baseball player.\nThe general purpose of the present invention, which will be described subsequently in greater detail, is to provide a new training bat system that has many of the advantages of the baseball training devices mentioned heretofore and many novel features that result in a new training bat system which is not anticipated, rendered obvious, suggested, or even implied by any of the prior art baseball training devices, either alone or in any combination thereof.\nTo attain this, the present invention generally comprises a tubular member having a center bore, a plurality of weight members removably positioned within the center bore, an inner cap secured to an inner end of the tubular member, and an outer cap secured to the outer end of the tubular member. A compression spring is preferably positioned between the weight members and the inner cap for maintaining the weight members non-movably adjacent one another. The tubular member is comprised of a first section having an outer diameter similar to a handle gripping of a baseball bat, a second section having a tapered structure, and a third section having an outer diameter smaller than said first section.\nThere has thus been outlined, rather broadly, the more important features of the invention in order that the detailed description thereof may be better understood, and in order that the present contribution to the art may be better appreciated. There are additional features of the invention that will be described hereinafter and that will form the subject matter of the claims appended hereto.\nIn this respect, before explaining at least one embodiment of the invention in detail, it is to be understood that the invention is not limited in its application to the details of construction and to the arrangements of the components set forth in the following description or illustrated in the drawings. The invention is capable of other embodiments and of being practiced and carried out in various ways. Also, it is to be understood that the phraseology and terminology employed herein are for the purpose of the description and should not be regarded as limiting.\nA primary object of the present invention is to provide a training bat system that will overcome the shortcomings of the prior art devices.\nA second object is to provide a training bat system for increasing the batting skills of a baseball player.\nAnother object is to provide a training bat system that increases a baseball player\"\"s mental and physical focus for making contact with a baseball.\nAn additional object is to provide a training bat system that may be utilized within various sports that utilize a bat to engage a ball such as but not limited to baseball, softball and similar sports.\nA further object is to provide a training bat system that improves the hand and eye coordination of a player.\nAnother object is to provide a training bat system that may be utilized by individuals of various ages, sizes and skill levels.\nOther objects and advantages of the present invention will become obvious to the reader and it is intended that these objects and advantages are within the scope of the present invention."}
{"text": "Liquid-liquid extraction (LLE) is a crucial step in the manufacturing of a wide range of products. Such processes are used for the extraction of one compound and as well as for separation between two or more compounds (fractional extraction).\nThe LLE process is simple in concept and usually requires the contacting of a feed containing the solute to be extracted with a solvent. This solvent/feed mixture is usually immiscible, but may be partially miscible.\nThe extraction and the stripping involve liquid-liquid contacting in which the droplets of one phase are initially dispersed in a second phase to facilitate mass transfer across the liquid-liquid boundary. Basically, there are two types of LLE units, those in which each individual stage is a separate unit termed \"mixer settlers\" and those in which several stages are integrated into one column. Multi-stage columns can be simple spray or packed columns, or can have stages equipped with various types of mixing devices separated by coalescence sections. The stage efficiency and the throughput of such devices are, obviously, directly related to the mass transfer and the coalescence rates.\nTo form small drops and ensure good contact between the phases, in slow mass-transfer systems, high intensity mixing is required. However, the shear stress induced by such a mixing can, in many cases, damage high molecular weight molecules. In addition, the intense mixing forms fine dispersions which reduce the coalescence rate or, in the presence of surface active impurities, may even cause a \"stable emulsion\", one of the operating hazards of solvent extraction equipment.\nFrom the point of view of LLE processes, the stability of the dispersion is its most important property, since the phases must separate at each extraction stage. For all practical purposes, the breakup time or the coalescence rate will determine the workable throughput of the extraction equipment. In countercurrent column-type contactors, steady operation is possible only when the rate of droplets arrival does not exceed the coalescence rate at the main interface; otherwise the dispersed band will extend over the entire column, leading to flooding. In mixer-settler contactors, the dimensions of the settler are designed according to the coalescence rate, and increasing the throughput above the coalescence rate will result in flooding of the settler. It is clear, therefore, that systems with emulsification tendencies cannot be operated by conventional extractors. Commonly such systems are handled in centrifugal extractors or by filtration. In some cases, adding compounds that break the emulsion can minimize the problem. This, however, makes complex the final purification of the product.\nEmulsion formation is a common problem in the pharmaceutical industry, where the desired products are frequently extracted from the fermentation broth by organic solvents and in extraction processes where mechanical agitation is used to increase the mass transfer rates.\nU.S. Pat. No. 4,954,260 of Z. Ludmer, R. Shinnar, and V. Yakhot (\"the U.S. Pat. No. '260 patent\") describes a process that overcomes some of these difficulties. In that case, special solvents are used that at one temperature form a homogeneous, one-phase mixture and, at a higher or lower temperature, form two phases, one solvent-rich and the other water-rich. See, also, Ullmann, Ludmer and Shinnar, \"Novel Separation Process Using Solvents with a Critical Point of Miscibility,\" Proc. A.I.Ch.E. Conference, Miami, Fla. (1993).\nThe process is composed of two stages, the heating Stage and the cooling stage where the mixture is cooled across the coexistence curve. Only very mild agitation is required. Compared to the mixing stage of the conventional isothermal extraction process, the heating stage provides a greatly improved contacting area, since in such stage there is but a single phase. Even more importantly, in the process of the U.S. Pat. No. '260 patent, the cooling stage has great advantage over the settling stage of the conventional extraction process: namely, when the cooling is fast enough, rapid phase-separation occurs even in the presence of impurities and cell debris. Accordingly, the need to use centrifuges is eliminated.\nUnfortunately, as a practical matter the process has the disadvantage of requiring continual and rapid heating and cooling."}
{"text": "The technical field of this invention is powertrain torque control for a motor vehicle.\nWhen a driven wheel of a vehicle slips with respect to the road surface, either by slowing down or speeding up, lateral adhesion of the tire to the road can decrease quickly and significantly. Loss of lateral adhesion can allow the tire to slip sideways and thus cause understeer (if a front wheel) or oversteer (if a rear wheel). Wheel slip is controlled in many vehicles by traction control systems during powertrain produced acceleration and by anti-lock braking controls during application of the vehicle brakes. But such slip may sometimes be produced by engine braking during vehicle deceleration when the throttle is closed with no vehicle brakes applied, when the engine slows down faster than the vehicle body and causes a braking torque to be applied to the driven wheels. Anti-lock braking controls are of no use when the vehicle brakes are not applied; and most acceleration traction controls are generally not designed to deal with wheel slip due to engine braking.\nOne system for sensing a difference in wheel speeds between driven and undriven wheels due to engine braking and increasing fuel to provide increased engine torque and spin up the slipping wheel has been suggested in U.S. Pat. No. 3,802,528; but modern computerized powertrain controls permit a far more finely tuned and accurate closed loop control to more advantageously balance the opposing goals of controlling wheel slip and providing engine braking.\nThe powertrain torque control of this invention provides an engine drag control mode of operation during periods of undesired engine drag induced wheel slip by modifying the torque of the vehicle engine in closed loop control to maintain a driven wheel speed at a predetermined target velocity lower than the vehicle speed by a target velocity difference providing as much engine braking as is consistent with a desired degree of lateral traction. The control derives a velocity error as the difference between the driven wheel speed and the target velocity and derives and delivers to the powertrain a torque command for reducing the velocity error.\nThe torque control determines the engine drag control mode in response to the wheel speed sensors, preferably causing entry of the engine drag control mode when the driven wheel speed that is closest to the vehicle speed falls below the target velocity while powertrain delivered torque and throttle position are below predetermined values indicative of deceleration. The driven wheel speed closest to the vehicle speed is preferably chosen because only one driven wheel is necessary for lateral traction and this wheel requires the smallest reduction in engine braking. The target velocity difference is preferably determined as a weighted difference between vehicle speed and vehicle turn curvature, the difference being reduced by the latter for quicker response in vehicle turns."}
{"text": "The invention relates to a circuit forming an active RC filter to be used as a band-stop filter in the high and very-high frequency domains, having a second order transfer function, comprising a circuit that includes a building block for performing a filtering and an amplifying function, and a building block for performing the summation V.sub.S of the circuit output signal and the signal V.sub.E applied to the input of the circuit.\nThe invention is applied to the realisation of integrated circuits forming filters, in the high and very-high frequency domains, that can be used, for example in one embodiment, as band-stop filters in frequency doublers for rejecting the unwanted signal at the fundamental frequency or, in another embodiment, as a band-stop filter having differential inputs and outputs; or, in yet another embodiment, be particularly suitable for a fine adjustment of the rejection.\nAn active all-pass filter is already known from the prior art from the publication by J. TOW, in \"IEEE Spectrum, December 1969, pp. 103-107\", entitled \"A Step-by-Step Active-filter Design\".\nThis publication describes, among other things, an all-pass filter realisation having a second order transfer function, which filter is shown in FIG. 8 of the said document. This circuit comprises three operational amplifiers arranged in series and is looped back. A fourth operational amplifier performs the summation of the output signal of the first operational amplifier and the input signal of the filter circuit\nThe operational amplifiers that are used for constituting this circuit have a gain which is most certainly infinite at low frequencies, but which becomes very low at high and very-high frequencies. Therefore, the circuit known from the said document has the disadvantage of not being suitable for use in these frequency domains.\nFurthermore, the looped back circuit has the disadvantage to cause the circuit to oscillate in certain conditions.\nFIG. 6 of the said document also shows an active band-pass filter that has similar characteristic features and thus exactly the same disadvantages.\nThe circuits known from the said document further have the following disadvantages: on the one hand, they are formed by a large number of transistors, which:\ncauses manufacturing to be costly in many cases\nrequires a large surface and is disadvantageous for use in integrated circuits,\nentails high consumption and is also disadvantageous for the use as mentioned above.\nOn the other hand, their characteristic frequency is not adjustable.\nFinally, the capacitors used in this circuit have considerable dimensions and render this circuit not integrable\nTherefore, the present invention has for its object to provide an active filter circuit that allows to get rid of these disadvantages and that specifically:\ncan operate at high or very-high frequencies;\nand in the case of the use at very-high frequencies is easy to integrate, requires a small surface and has little consumption;\nexhibits characteristic features so that the rejection frequency and the rejection are adjustable and easy to control;\ncan admit a differential input signal and a differential output signal."}
{"text": "In an electronic system, for example, a solar photovoltaic assembly, communication assembly, etc., for transmitting a micro electric signal, it is necessary to transmit the micro electric signal to an electrical device disposed inside a connection box through a conductive sheet (bus bar) so as to collect or retransmit the micro electric signal. Inside the connection box, a conductive terminal for electrically connecting the conductive sheet and a connection piece connected between the conductive terminal and the electrical device are disposed.\nWith respect to FIG. 1, a conventional connection box 200 is shown and includes a housing 201, a conductive terminal 203 and a conductive sheet 204 (i.e. bus bar). A cable receiving passageway is formed in the housing 201 to introduce an external cable 202 therein. The conductive terminal 203 is electrically connected to the introduced cable 202 by, for example, soldering. The conductive sheet 204 is connected between the conductive terminal 203 and an electrical device (not shown), for example, a solar photovoltaic assembly. The conductive sheet 204 is made of a thin metal material, for example, a copper foil, and is soldered to the conductive terminal 203. Since the conductive sheet 204 is very thin, there is a possibility that a rosin joint occurs between the conductive sheet 204 and the conductive terminal 203 during soldering them, which may cause a poor electrical contact and an unstable electrical connection between them.\nFurthermore, it may produce an over high contact impendence and even an electric arc between them, which creates a serious potential safety risk. Also, soldering the conductive sheet 204 to the conductive terminal 203 is very complicated and difficult, and it is difficult for them to be separated after the soldering, which is unfavorable for maintenance of the solar panel of the solar photovoltaic assembly."}
{"text": "In communications systems, in particular in wireless radio systems, co-channel interference (CCI) and carrier frequency offset (CFO) are well-known effects which degrade the quality of the received signal. CCI is caused by an interferer signal having substantially the same carrier frequency as a user signal to be detected. Typically, CCI suppression is accomplished by filtering the received signal with an appropriate filter structure.\nCFO can be present in the user signal or the interferer signal or in both signals. CFO causes CCI suppression to deteriorate quickly with increasing distance from a training-sequence in the transmitted signal, leading to a decrease in signal quality. Typically, CFO compensation is done by measuring CFO and correcting the received signal by the measured frequency offset. CFO compensation is used for carrier frequency tracking.\nThe efficiency of CFO compensation depends on the CFO measurement equipment and on the algorithm and mathematical procedure for calculating the CFO correction quantity to be applied for signal frequency correction."}
{"text": "Wireless communication can be used as a means of accessing a communication network. Wireless communication has certain advantages over wired communications for accessing a network. For example, implementing a wireless interface can eliminate a need for a wired infrastructure thereby reducing the cost of building and maintaining network infrastructure. In addition, a wireless network can support added mobility by allowing a wireless device to access the network from various locations or addresses. A wireless interface can comprise at least one transceiver in active communication with another transceiver that is connected to the network.\nVarious types of network configurations can be used to communicate data over the wireless network. For example, a heterogeneous network can be configured to include various types of access nodes such as a macro access node, a micro access node, a pico access node, a femto access node, etc. In a heterogeneous network, a wireless device can be served by an access node having the lowest signal path loss rather than by an access node having the strongest signal strength as in traditional network configurations.\nIn a heterogeneous network, interference can occur at the cell edge of the short range, low power access nodes due to the macro access node. This interference can result in undesirable reduction in coverage and throughput to the wireless devices in communication with the short range access node. A scheduling scheme comprising almost blank subframes (ABS) can be used to create an opportunity for the wireless devices within the cell edge region of a short range access node to receive downlink information without interference from the macro access node. However, ABS subframes can undesirably limit an amount of resources allocated to wireless devices during each frame."}
{"text": "Game animals, particularly those hunted for food and/or sport, such as but not limited to deer, are attracted by scents varying from those of other animals and odorous attractant scents such as those of food. It is these scents which hunters use to pull animals closer for harvest opportunities.\nDuring mating seasons, male animals, as part of the mating ritual, attempt to attract females by “scraping” the ground, using their hoofs, at desirable locations and urinating in the scrape in an attempt to attract females. The females in turn, when attracted, deposit a female hormone in the “scrape” which is highly attractant to males. Hunters, in an attempt to mimic these attractants, have developed commercially available “attractant scents” which substantially duplicate the male and female mating scents.\nThose who wish to examine animals at close range, such as hunters, routinely manually disperse such scents on the ground in an attempt to attract their quarry.\nIn some cases, hunters prepare the ground, by making a mock scrape using various implements, to simulate a deer “scrape” before dispensing attractant scent liquid into the scrape. However, humans entering these sensitive areas can disturb the animals, in ways known only to animals. In many cases, these trespasses into an animal's area provide persistent unseen warnings which tend to keep the desired animal from approaching.\nVarious devices have been developed over the years which attempt to continuously or periodically deliver portions of scented liquids to the desired spot without human interaction. The present invention provides significant improvements over the prior system as described below.\nU.S. Pat. Nos. 5,279,062; 5,361,527 and 5,220,741 all to Burgeson describe devices for dispensing a scented liquid (scent) onto the ground. The device employs a rigid camouflaged scent container having a cap with an exterior nozzle tube which may be straight or bent through 180 degrees into a J shape or through 360 degrees into a circular shape. The container is suspended over the ground and is partially filled with the scented liquid. As the air in the space over the liquid expands during the day it pushes out a volume of scented liquid. However, these devices have several problems.\nThe Burgeson patents suffer from the defects that the scent container can be only partly filled. Also it must be clear that the amount of liquid scent delivered depends on the unfilled volume within the container. When the container is more filled, the air volume remaining is less. Thus, less liquid scent can be delivered for a given temperature change. Principle operations of these devices are dependent on temperature changes.\nAdditionally, the Burgeson devices require a rigid container must be used, when the container is filled, it has a large mass and must respond more slowly to any temperature change. By contrast, when the container is nearly empty, there is a large gas volume within the container which will cause a larger amount of liquid scent to be delivered for a given temperature change. Further, when the container is nearly empty, the small mass of liquid heats easily, thereby causing a widely varying rate and quantity of liquid scent to be delivered, depending on the fraction of the bottle which is filled.\nFurthermore, the containers in the Burgeson patents must be made of some rigid material such as glass. Such a container is easily susceptible to being broken, should the container fall from the tree or support where it is suspended. Additionally, if the container is made of rigid plastic, the container can crack over time with constant exposure to sunlight or other environment factors, such as, heat, cold, or changes in temperature over time.\nU.S. Pat. No. 8,510,984 to Burgeson describes a temperature regulated, pressure activated liquid scent dispenser. The pressure in the interior of the container increases as ambient temperature increases. A release structure of the container will release a portion of the liquid scent once a threshold pressure or threshold amount of pressure build-up is reached in the interior of the container.\nThis device requires filling the interior volume with a liquid scent so that the interior volume also includes a volume of air, suspending the dispenser over a ground surface, and dispensing the liquid scent from the interior volume through the release structure. Due to an increase of pressure of the volume of air, and upon reaching a threshold air pressure, the release structure releases a portion of the liquid scent from the interior volume.\nSimilar to the other Burgeson Patents: U.S. Pat. Nos. 5,279,062; 5,361,527 and 5,220,741 these devices provide for liquid delivery dependent on temperatures changes to shift atmospheric pressure inside the rigid vessels in order to drive the liquid dispersal.\nU.S. Pat. No. 8,739,455 to Burgeson describes a temperature regulated, pressure activated liquid scent dispenser. The pressure in the interior of the container can increase as ambient temperature increases. A release structure of the container releases a portion of the liquid scent once a threshold pressure or threshold amount of pressure build-up is reached in the interior of the container.\nThe Burgeson '455 patent device provides for filling the interior volume with a liquid scent so that the interior volume also includes a volume of air, suspending the dispenser over a ground surface, and dispensing the liquid scent from the interior volume through the release structure. Due to an increase of pressure of the volume of air, and upon reaching a threshold air pressure, the release structure releases a portion of the liquid scent from the interior volume.\nSimilar to the previously referenced Burgeson, this device includes a, pressure interior to the vessel that is dependent on ambient temperature variations to create internal vessel pressure to drive the liquid delivery system.\nU.S. Pat. No. 5,971,208 to Kennedy describes a device for delivering an animal attractant scented liquid employing a flexible walled container with an external gas filled balloon strapped to the container so positioned that expansion and contraction of the gas within the balloon, in response to temperature changes, causes the wall of the container to flex so as to discharge liquid from the container when the temperature rises, and cease discharging on a temperature drop.\nSimilar to the other prior art is in its use of atmospheric pressure and temperature changes to expand the gas inside the dispersal drive mechanism. Here, the Kennedy '208 patent uses externally mounted balloons.\nU.S. Published Patent Application. 20080054021A1 to Brown et al. describes a product directed at the deer hunting market for scent application and dispersal that uses a molded, rigid container capable of being filled with liquid scent and then dispensed in a multiple ways, including a flip top cap for direct placement of its contents to a variety of specific areas, using scent wicks, cotton balls, etc. to establish scent stations (dipping these into open liquid reservoir). Similar to the other prior art, Brown '402 relies on temperature changes to drive fluid dispersal.\nThus, the need exists for solutions to the above problems with the prior art that does not require head space nor temperature changes to drive animal attractant liquid therefrom."}
{"text": "1. Field of the Invention\nThis invention relates to dental implements, and more particularly, to a toothpick holder for holding a toothpick in any one of several different adjusted positions for easier access to difficult to reach areas of the mouth.\nMore specifically, the present invention relates to a toothpick holder which includes a handle with a toothpick retainer projecting perpendicularly therefrom and rotatable about its axis to several different latched positions. The retainer has a pair of aligned openings therein for snugly receiving a toothpick of yieldable material, whereby the toothpick is held in a firmly latched position with its axis extending in a desired direction relative to the axis of the handle. By manipulating a latch member, the retainer may be released for removal from the handle or adjustment to different latched positions.\nIt is well known by the dental profession that brushing does not always adequately clean the teeth, particularly in the areas between the teeth. Thus, flossing and other cleaning methods are recommended in conjunction with brushing. Moreover, the proper use of toothpicks is very beneficial in any oral hygiene program, and can be particularly effective in cleaning the spaces between the teeth.\nHowever, except for a few attempts at developing a toothpick holder, people are generally limited to the use of wood or plastic toothpicks held in the users hand. Accordingly, the use of a toothpick is only partially effective in cleaning the teeth, and those areas which are difficult to reach are usually not cleaned.\nPrior Art\nExamples of prior art toothpick holders are shown in U.S. Pat. Nos. 710,498, 1,291,282 and 3,892,040. In U.S. Pat. No. 710,498 a quill-like member is inserted through a shaped holder whereupon the quill-like member is curved to form a pick. U.S. Pat. No. 1,291,282 discloses a threaded holder having a pair of openings therein for receiving a toothpick in either of two different positions. U.S. Pat. No. 3,892,040 discloses a holder having a threaded sleeve which is movable against a round toothpick to clamp the toothpick in position.\nU.S. Pat. No. 3,471,929 discloses a dental implement in which a shaft 26 is held to a handle by a pin 8. A blade 30 is carried by the shaft for performing gum cutting operations.\nNone of the above patents teaches a toothpick holder capable of holding a toothpick with a wedge shaped cross section in any one of several different adjusted latched positions, and with the particular cooperation between the elements and pick as set forth more fully hereinafter, wherein the holder properly orients such a wedge shaped toothpick for optimum effectiveness in all areas of the mouth."}
{"text": "Swimming pools, both in ground and above ground, are becoming more prevalent in sections of the country where freezing occurs during winter months. Conduits and equipment connected to the pools, for examples drain lines, recirculating lines and the like, must be drained of all water prior to the freezing season to insure that the conduits and equipment will not freeze and burst.\nThe most popular prior method of winterizing a swimming pool was to lower the level of the water within the pool to an elevation below the drain line. Often, the water in the drain line was then drained back into the pool or in the alternative removed at, for example, the filtering equipment.\nSeveral thousand gallons of water were lost every year by the lowering of the water levels in the swimming pools to the required elevation to drain the conduits. Besides being time consuming, the cost of water has increased and the overall loss of water was very expensive. In addition, because this water was treated water, its loss was not only expensive but detrimental from a conservation point of view."}
{"text": "When a connecting tube of a blood pump or a humor circulating circuit and a tube such as a cannula mounted on a circulatory system organ are connected in order to pour blood or humor substitution fluid in the circulatory system organ, it is necessary to carefully connect the tubes so as not to allow bubbles to remain within the tubes. If the bubbles are contained in blood, a blockade of blood capillaries or the like occurs, resulting in an extremely dangerous condition for a living body.\nAs one means for connecting two tubes filled with blood or the like to eliminate bubbles therein, there is a method for connecting the tubes while pouring the humor substitution fluid or physiological saline solution to connected ends of two tubes. In accordance with this method, a tube such as a cannula mounted on a blood vessel is supported with a connected end directed upwardly, blood is moved up to the interior of the tube by blood pressure, and thereafter, the tube is clamped below a liquid level of blood. Then, the humor substitution fluid is filled from the clamped portion to the upper connected end, and the tube is sufficiently commoved to completely discharge bubbles adhered to tube walls from the interior of the tube. Similarly, the connecting tube of the blood pump is also filled with the humor substituion fluid to completely discharge bubbles, two tubes are made to come closer to each other with connected ends thereof directed upwardly, and the connected end of one tube is inserted into the connected end of the other tube within the flowing-down humor substitution fluid while a large quantity of the humor substitution fluid to said portions which are formed to come closer to each other to connect two tubes.\nIn the above-described connecting method, even if bubbles are carefully discharged by commotion, air tends to be entrained into liquids since connecting of tubes is carried out within the flowing-down liquid and a part of the air is drawn into the tubes and formed into bubbles. Thus, there was a disadvantage in that the bubbles remain within the tubes despite the fact that connecting of tubes was carried out with effort.\nIn view of the foregoing, a connecting method has been used in which a connecting tube provided with branch tubes is used to connect two tubes, and bubbles within the tubes after connection are guided into the branch tubes for discharge. In this method, a cannula from which bubbles within the liquid have been discharged and a connecting tube are connected by insertion with the branch tubes directed upwardly, the connecting tube is curved in the form of a mountain to bring the branch tubes to the highest position, the two connected tubes are commoved to upwardly move the bubbles up to the branch tubes to completely remove the bubbles within the connected tubes, and thereafter, root portions of the branch tubes are clamped.\nThis method making use of the branch tubes carries out the discharge of bubbles after the tubes have been connected and therefore can positively prevent the bubbles from being retained as compared with the method of connecting tubes after removal of bubbles. However, vestiges remain on side walls of portions where the branch tubes were present and thrombus tends to occur on said vestiges. Therefore, there involves a disadvantage that in case of transfusion of blood, heparin has to be contained in blood.\nAs the result of a series of studies on means and apparatus for connecting two tubes which can positively remove bubbles within the tubes and which can prevent re-generation of bubbles when the tubes are connected and remain on vestiges on side walls of the tubes connected, the present inventors have found that a method can be used in which two tubes are connected within the liquid in order to facilitate connection of the tubes with bubbles being eliminated completely and positively."}
{"text": "Chlorine dioxide, ClO2, is one of the most effective bleaching agents for use in industrial and domestic process and services, and for commercial and consumer products. The strong oxidative potential of the molecule makes it ideal for a wide variety of uses that include disinfecting, sterilizing, and bleaching. Concentrations of chlorine dioxide in an aqueous solution as low as 1 part per million (ppm) or less, are known to kill a wide variety of microorganisms, including bacteria, viruses, molds, fungi, and spores. Higher concentrations of chlorine dioxide, up to several hundred ppms, provide even higher disinfection, bleaching and oxidation of numerous compounds for a variety of applications, including the paper and pulp industry, waste water treatment, industrial water treatment (e.g. cooling water), fruit-vegetable disinfection, oil industry treatment of sulfites, textile industry, and medical waste treatment.\nChlorine dioxide offers advantages over other commonly used bleaching materials, such as hypochlorite and chlorine. Chlorine dioxide can react with and break down phenolic compounds, and thereby removing phenolic-based tastes and odors from water. Chlorine dioxide is also used in treating drinking water and wastewater to eliminate cyanides, sulfides, aldehydes and mercaptans. The oxidation capacity of ClO2, in terms of available chlorine, is 2.5 times that of chlorine. Also, unlike chlorine/hypochlorite, the bactericidal efficiency of chlorine dioxide remains generally effective at pH levels of 6 to 10. Additionally, chlorine dioxide can inactivate C. parvum oocysts in water while chlorine/hypochlorite cannot. Hypochlorite and chlorine both react with the bleached target by inserting the chlorine molecule into the structure of the target. Though this mode of reaction can be effective, it can result in the formation of one or more chlorinated products, or by-products, which can be undesirable both from a economic sense (to eliminate hydrocarbons from the reaction media) and a safety and environmental standpoint. In addition, the step of bleaching by hypochlorite and chlorine results in the destruction of the bleach species itself, such that subsequent bleaching requires a fresh supply of the chlorine bleach. Another disadvantage is that certain microorganisms that are intended to be killed by these two commonly-used bleach materials can develop a resistance over time, specifically at lower concentrations of the chlorine or hypochlorite.\nChloride dioxide is generally used in an aqueous solution at levels up to about 35%. It is a troublesome material to transport and handle at high aqueous concentrations, due to its low stability and high corrosivity. This has required end users to generate chlorine dioxide on demand, usually employing a precursor such as sodium chlorite (NaClO2) or sodium chlorate (NaClO3).\nA typical process for generating chlorine dioxide from sodium chlorate salt is the acid-catalyzed reaction:NaClO3+2HCl→NaCl+½Cl2+ClO2+H2O\nSodium chlorite is easier to convert to chlorine dioxide. A typical process for generating chlorine dioxide from sodium chlorite salt is the acid-catalyzed reaction:5NaClO2+4HCl→4ClO2+5NaCl+2H2O\nFurther details on the acid-catalyzed reactions of chlorites and chlorates to produce chlorine dioxide can be found in “Chlorine Dioxide Generation Chemistry” (A. R. Pitochelli, Rio Linda Chemical Company), Third International Symposium: Chlorine Dioxide Drinking Water, Process Water and Wastewater Issues, Sep. 14, 15, 1995, La Meridian Hotel, New Orleans, La., incorporated herein by reference.\nA common method of making chlorine dioxide uses a multi-chamber electrolysis cell that converts the chlorite salt into chlorine dioxide. This method uses separately an anode compartment and a cathode compartment that are separated by an ion permeable membrane. The separate compartments operate with significantly different reactants, and contain solutions with different pH values. One example of a multi-compartment electrolysis cell is disclosed in U.S. Pat. No. 4,456,510, issued to Murakami et al. on Jun. 26, 1984, which teaches a process for forming chlorine dioxide by electrolyzing a solution of sodium chlorite in an electrolysis cell that contains an anode compartment and a cathode compartment separated by a diaphragm, preferably a cation exchange membrane. Another example of a two-chamber electrolysis cell is disclosed in U.S. Pat. No. 5,158,658, issued to Cawlfield, et al. on Oct. 27, 1992 which describes a continuous electrochemical process and an electrolytic cell having an anode chamber having a porous flow-through anode, a cathode chamber, and a membrane there between.\nWhile separate-compartment, membrane-containing electrolysis cells have been used to make chlorine dioxide on a commercial scale, they have not been completely satisfactory. Even though they may have convenience advantages over the conventional acid catalysis production of chlorine dioxide, the electrochemical approach has proven to be more expensive to produce large volumes of chlorine dioxide. The electrolysis cells in commercial use, and disclosed in the prior art that utilize ion permeable membranes or diaphragms, require that the anolyte solution be substantially free of divalent cations, such as magnesium and calcium, to avoid the formation of precipitated calcium or magnesium salts that would quickly block and cover the membrane, and significantly reduce or stop the electrolysis reaction.\nConsequently, there remains a need for a simple, safe method and apparatus for manufacturing chlorine dioxide to meet a wide variety of commercial and domestic uses, under a wide variety of situations. The present invention describes a method and an apparatus for making chlorine dioxide inexpensively, easily and effectively."}
{"text": "Bed-type massage devices are generally provided with massaging rollers which can be displaced by a drive device in a bed base and are so constructed that the massaging rollers travel on guide rails provided on both sides inside the bed base, the arrangement being such that when a user lies face upwards on the bed-type massage unit the massaging rollers roll along his or her dorsal region and effect finger-pressure type massage thereof.\nIn recent years there has been research on and development of units in which such bed-type massage devices are made contractible to as compact a size as possible in order to reduce the space needed at times of transport and delivery or during storage.\nFor example, as disclosed in Japanese Laid-open Patent Application No. 59-189848 and Japanese Laid-open Utility Model Application No. 61-54834, there are known bed-type massage devices which are designed to reduce space requirements at times of transport and delivery or during storage by being constructed in such a way that foldable guide rails are provided on both sides inside a bed base and the bed frame itself is formed as an elastic body which is foldable and the massaging rollers can run along the guide rail, whereby the base unit can be folded into halves or thirds when it is not in use.\nHowever, the structure in these conventional bed-type massage devices is such that extension and contraction of the bed base is employed solely for the purposes of transport and delivery or storage and, although the size of the bed base can be reduced in the longitudinal direction, it is impossible to reduce the volume and so the devices still fail to resolve problems in packing, etc.\nFurther, since conventional bed-type massage devices are units designed for the purpose of reduction of size at times of transport and delivery or during storage and are not constructed in a manner permitting adjustment of the extension or contraction in accordance with the height of the user, if a user who is shorter than the length of the bed base uses the base to effect massaging, the massaging rollers move to portions beyond the body of the user, i.e., they move over an unnecessarily large range, which is wasteful both in terms of electric power and of time.\nBy way of a means for resolving this problem, the present Applicant discovered a means whereby a bed base of a bed-type massage device can be extended or contracted in opposed lengthwise directions without being folded and which permits fine adjustment of the extension in accordance with the height of the user. In this case, however, although extension and contraction in opposed lengthwise directions and fine adjustment of the extension are possible, there are problems associated with aspects such as the means for disposing the massaging rollers inside the bed base to match the bed base in different states and the drive means for causing the massaging rollers inside the bed base to move forward and back in a manner such as to match the bed base in different states.\nIt is the object of the present invention to resolve the various problems noted above and provide a bed-type massage device in which a massage unit is not driven when the bed base is contracted but the massage unit can longitudinally travel repeatedly and smoothly over the whole area of the bed base when the bed base has been extended or contracted and fine adjustment of the extension has been made.\nIt is a characteristic of the bed-type massage device of the invention that it comprises a variable bed base which is made freely extendible and contractible in opposed lengthwise directions and in which a first side of a flat second base fits into an opening on one side of a flat first base, at least one retention hole is formed on the left and on the right in the lower surface of the second base, and lock pins which are engageable in these retention holes are so provided that lock mechanisms provided at the tips thereof face the second base from the left and right of the other side of the first base; a pair of guide rails which are laid lying along the longitudinal direction on the left and right of the upper surface of at least the first base of the said bed base; a drive mechanism in which a drive motor is provided at and orthogonal to one end of the second base in the longitudinal direction, worms are respectively coupled with the ends of two drive shafts of this drive motor and the respective ends of a pair of rod-like screwshafts which are respectively provided along both lengthwise sides of the second base are in screw engagement which the respective worms; a massaging unit which is constituted by effecting screw engagement of both sides of the base end of a frame unit with the pair of rod-like screwshafts of the said drive mechanism providing facing running rollers which are capable of moving and running over the said guide rail on both sides of the lower part of the far end of the said frame unit and providing rolling members at a set interval above the base end and the far end of the said frame unit; and a covering cloth that is wound around to cover the entire upper surface of the bed base, and the device is so constructed that the massage unit can repeatedly travel over the bed base and the variable bed base is freely extendible and contractible to permit adjustment to any required length.\nHaving the above construction, the bed-type massage device of the invention brings about the following effects.\nAll that is needed if it is wished to adjust the length of the bed-type massage device of the invention to match the height of a user is to push the first base of the variable bed base in the direction of the second base and lock it in the said position by means of the lock mechanisms, and if it is required to shorten the bed at times such as when the bed is packed for transport or is delivered it can easily be reduced to about half its length by releasing the lock mechanisms and pushing in the first base into which the second base is inserted.\nFurther, when it is wished to extend the bed base after it has been shortened, this can easily be done by simply releasing the lock mechanisms, pulling the first base so that the second base is moved out from it and then, after the bed has been set to the required length, locking it with the lock mechanisms.\nWhen the adjustment of the length of the bed-type massage device has been completed, the user lies face upwards on its variable bed base and then simply actuating the drive motor brings about an agreeable massaging action in which, through the action of the drive mechanism, the massage unit is caused to move and travel on the guide rails, limit switches at the variable bed base's terminal and start ends in the longitudinal direction are actuated and cause repeated displacement and travel of the said massage unit, so causing rotatable rollers of the rolling members to slide while coming into uniform contact with the whole surface of the user's back.\nThe bed-type massage device of the invention will now be described in detail with reference to one embodiment thereof which is shown in the drawings."}
{"text": "The invention concerns a device for setting the ignition timing in projectile fuses. The device comprises a timing ring rotatable around the front-to-rear axis of the fuse for setting the ignition time.\nFrom German Pat. No. 22 00 540, a device for setting the timing of a projectile fuse is known. The disclosure of that patent is incorporated herein by reference. The device consists of a stationary point, a rotatable part and a stationary lower part. The rotatable part comprises a timing ring and is used to get a predetermined ignition time or ignition delay. The timing ring may be rotated by a maximum of 360.degree. and is equipped with a device whereby it is possible to set even longer time delays with a full rotation. To prevent the unintentional resetting of the set ignition time or delay time during subsequent manipulation of the fuse, special measures are usually effected by means of threaded bolts located on a circular ring.\nA disadvantage of this arrangement involves the fact that, firstly, the assembly is relatively complex. Secondly, a nonuniform screwing-in of the bolts cannot be entirely prevented; the axial forces are fully applied to the individual threaded bolts and therefore deformations on the housing cannot be entirely avoided.\nIt is an object of the invention to provide a device of the afore-mentioned type, whereby the time setting torque customarily used in time fuses may be set exactly, and additionally a positive connection of the center and lower parts of the housing is possible."}
{"text": "The present invention is directed, in general, to a method for guiding the drilling of wells at a substantial depth in the earth, and more particularly to methods for determining the distance and direction to a target-well from a borehole being drilled.\nThe difficulties encountered in guiding the drilling of a borehole to intersect, to avoid, or to parallel an existing well at distances of thousands of feet below the surface of the earth are well known. Such guidance may be required when it is desired to avoid existing wells in a field, or when existing oil or gas wells have blown out and it becomes necessary to drill intersecting relief boreholes to prevent serious damage to underground gas or oil fields. Various electromagnetic methods for the precise drilling of such relief boreholes have been developed and have met with significant success during the past few years. Such methods and the instruments used are described, for example, in U.S. Pat. No. 4,323,848 and U.S. Pat. No. 4,372,398, both issued to the applicant herein, and in U.S. Pat. No. 4,072,200 to Morris et al. See, also, Canadian Patent No. 1,269,710 of Barnett et al, issued May 29, 1990.\nEven though the guidance of boreholes with respect to existing wells is, in general, well developed, special problems can occur where existing techniques are not sufficient to provide the precise control required. For example, when it is desired to locate and to either avoid or to intersect a particular target well in a field which includes numerous other wells, problems can occur. Thus, when multiple wells lead from a single location, such as a drilling platform, it may become necessary to drill a borehole to avoid intersecting neighboring wells or, alternatively, to intersect a particular one of, for example, sixteen wells, all starting at approximately the same location and spreading downwardly and outwardly from each other. The borehole being drilled may start at the same general location as the other wellheads, or may start at a location several hundred feet from the wellhead of a target well. If intersection with a specific well is desired at, for example, three thousand feet below the surface, guidance information can be provided by a low-frequency alternating current injected into the earth, as from an electrode in the borehole being drilled, with the resulting earth current being concentrated in a casing or other electrically conductive material at the target well. The current so produced in the target well results in a magnetic field which can be detected from a highly sensitive magnetometer located in the borehole. However, in multiple-well fields, the use of such a current injection system results in a target current being induced in all of the wells in the region, not just the target well. This produces multiple magnetic fields which are superimposed at the borehole magnetometer, making it extremely difficult to obtain accurate distance measurements to the target well of interest, thus interfering with drill guidance.\nProblems are also encountered in drilling non-parallel wells, such as a horizontal well through a field of vertical wells, or vice versa, where it is desired to avoid the existing wells, or in the alternative to intersect a specific well. Another area of difficulty occurs in the drilling of multiple horizontal wells, particularly where a well being drilled must be essentially parallel to an existing well. The need to provide two or more horizontal wells in close proximity, but with a precisely controlled separation, occurs in a number of contexts such as in steam assisted recovery projects in the petroleum industry, where steam is to be injected in one horizontal well and mobilized viscous oil is to be recovered from the other. This process is described, for example, in Canadian Patent No. 1,304,287 of Edmunds et al, which issued Jun. 30, 1992. Another example is in the field of toxic waste reclamation, where there is a need to drill parallel horizontal wells under waste disposal sites so that air can be pumped into one and toxic fluids forced into and recovered from the other. Again, in hot rock geothermal energy systems, there is a need to drill parallel wells so that cold water can be injected into one and not recovered from the other.\nThe need to drill horizontal, parallel wells is of most immediate concern in the mobilization of heavy oil sands, where a borehole is to be drilled close to and parallel to an existing horizontal well with a separation of about ten meters for a horizontal extension of a thousand meters or more at depths of, for example, 500 to 1,500 meters. A number of such wells may be drilled relatively closely together, following the horizon of the oil producing sand, and such wells must be drilled economically, without the introduction of additional equipment and personnel."}
{"text": "1) Field of the Invention\nThe present invention relates to a timing analysis apparatus, a timing analysis method, a timing analysis program, and a recording medium for static timing analysis of circuits arranged and wired on a large-scale integration (LSI) chip.\n2) Description of the Related Art\nIn the conventional technology, in static timing analysis (STA) of a LSI, on chip variation (OCV) of delay time (or delay) in a LSI chip arising from temperature, voltage, and manufacturing process is used.\nIn concrete, a variation on this LSI chip is quantified, and is used as a variation value. For example, by multiplying delay time of a path to input to a sync flip-flop (FF), a path to input to a source FF, and a path between both the FFs by a variation value, timing analysis is carried out. It is necessary to make this variation value a value in which all the variations of supposed delay time are taken into consideration. Among them, with regard to manufacturing process variations, a distance between circuit elements that cause a variation value is set with a supposed maximum range as a reference. Meanwhile, with regard to voltage variations, a potential difference owing to a power source decline is set as a variation value. Furthermore, with regard to temperature variations, a temperature difference between circuits is set as a variation value.\nA timing analysis method according to the conventional technology is explained. FIG. 1 is a flowchart showing a timing analysis method according to the conventional technology. Macro cells and logic circuits are arranged and wired on a LSI chip (step S1101). From layout data obtained from this arrangement and wiring, wire resistance values and inter-wire capacity values are extracted (step S1102). Then, by use of the layout data, wire resistance values, and inter-wire capacity values, a delay time is calculated (step S1103). By use of this calculated delay time and the variation values mentioned above, timing analysis is carried out (step S1104). Finally, an error obtained by the result of timing analysis is corrected (step S1105).\nIn fully custom LSI designs, a delay characteristic analysis method is disclosed where prompt and precise delay characteristic analysis is enabled, when conditions concerning circuit designs and signal transmissions are changed partially, by making the most of analysis results on circuits before the change (see, for example, Patent Application Laid-Open Publication No. 2002-215710). Furthermore, a method to attain both a shortened development period and a low electricity consumption that are in relations of a trade-off is disclosed (see, for example, Patent Application Laid-Open Publication No. 2002-312410).\nHowever, in the methods according to the conventional technology, variation values used in timing analysis are values in which all the variations of supposed delay time, i.e., variations of voltage, temperature, and manufacturing process, should be taken into consideration, therefore, the values are used as constant ones to any of circuits.\nAccordingly, when the range of a circuit that causes variation values is narrow (for example, the distance between a sync FF and a source FF is short), variations of manufacturing process and voltage become excessive, and as a result a useless margin occurs. The useless margin makes it difficult to converge the timing, and also increase the size of LSI chips, which has been a problem with the conventional technology."}
{"text": "1. Field of the Invention\nThe present invention is directed towards a protective cover dimensioned and configured to overlie the keyboard of a computer while enabling the user of the computer to have free access thereto. Pets such as cats or the like are thereby prohibited from inadvertently walking or jumping on the keys of the keyboard both during operation of the computer and during periods of non-use.\n2. Description of the Related Art\nIt is well understood among pet owners and particularly those individuals who are attracted to cats that in spite of the beneficial relationship existing between an owner and his pet, there are certain inconveniences or annoyances frequently associated with living and caring for such animals. Cats, particularly those of a curious or affectionate nature, frequently insist in occupying the same space or surrounding vicinity as the owner regardless of the activity in which the owner is involved. For example, it is quite common for a cat, of an affectionate nature to spring into the lap of the owner or caretaker while the individual is reading a newspaper or the like. Also, when not directly involved in the activities of the owner, cats in particular, frequently prowl their domain and in doing so interfere with household objects by either purposely or accidently displacing certain objects.\nWith the advent of the modern day personnel computer, the popularity thereof have increased to the point where almost a majority of the households and work places include some type of computer facility which includes a keyboard. This keyboard is normally operated from some type of table or horizontal work surface and whether the computer is being utilized or not the keyboard is generally left in an exposed position. These individuals who are pet owners and particularly cat lovers and also use some type of computer instrument involving a keyboard have encountered the problem of having cats generally prowl across the keyboard and/or position themselves firmly on the keyboard thereby disturbing the operation thereof. In addition, cats of a particularly curious nature or those who have the habit of invading the operating area while their owners are working on the computer certainly have a disruptive affect on the operation and in fact, could, in certain instances, destroy certain work currently being performed on the computer.\nAccordingly, there is a recognized need in this area to protect computer keyboards to the extent that pets, such as but not limited to cats, would be prevented from coming in contact with the keyboard. This protection should occur during periods when the computer is being used and the keyboard is being manipulated and also during periods of non-use. Moreover, a protector assembly should be of generally lightweight construction so as to be easily useable, even by children in the protection of the keyboard, while also making it substantially difficult for an animal to have access to the keyboard and/or remain in an obtrusive location relative to the keyboard. Further, such a device should also have sufficient stability such that investigation by a cat or other small animals would not serve to remove the protection, while allowing clear access to the keyboard by an individual in a manner that the user's hands are free to operate or manipulate the keyboard in the conventional fashion. Also, at least a portion of the overall assembly should be structured to allow clear, and unobstructed viewing of the entire operative surface of the keyboard during operation of the computer and utilization of the keyboard."}
{"text": "1. Field of the Invention\nThe present invention relates to a near-infrared shield and a display front plate using the near-infrared shield.\n2. Description of Related Art\nRecently, demands for plasma display panels (PDPs) as display panels for various electronic equipment like large-sized TVs have been increased. A PDP includes two glass plates on which a fluorescent substance is applied, and a gap between the glass plates is filled with a gaseous mixture containing xenon and neon. When a high voltage is applied to the gaseous mixture, ultraviolet radiation is generated. The ultraviolet radiation impinges the fluorescent substance, and thus the fluorescent substance emits light.\nAt this time however, in addition to the ultraviolet radiation, near-infrared radiation in a wavelength region ranging from 820 nm to 1100 nm, electromagnetic waves and the like are generated as well. Since the near-infrared radiation has a wavelength region overlapping the wavelength region used for near-infrared communications or remote controls of the other electronic equipment, it can cause malfunctions of the equipment. For solving the problem, a near-infrared shield is provided on the front plate of the PDP so as to absorb the near-infrared radiation (see “Characteristics of an antireflection film and optimum design/film formation technology” by Hanaoka et al., first edition, second printing at Technical Information Institute Co., Ltd., Feb. 5, 2002, p. 184).\nFor the near-infrared shield, for example, a product prepared by dispersing a near-infrared absorption compound in a resin and shaping it into a film is known. Examples of the near-infrared absorption compound include a diimonium compound, a phthalocyanine compound, a cyanine compound and the like. The compounds are known for exhibiting a particularly excellent near-infrared absorption characteristic in a case where two or more of the compounds are used in combination, particularly in a case of combining the diimonium compound and either the phthalocyanine compound or the cyanine compound in comparison with a case of using one of the compounds alone (see JP 11(1999)-316309 A and JP2003-21715 A).\nIn general however, conventionally-used phthalocyanine compounds do not have the desired solubility in solvents or compatibility with resin. Therefore, various substituents must be introduced when such a compound is used for a near-infrared shield, and this increases the production cost. The conventionally-used cyanine compound is easy to obtain. However, when it is combined in use with a diimonium compound, mutual interaction is generated between the two near-infrared absorption compounds in a long-term storage, and this will cause a problem that both the near-infrared radiation absorptivity and the visible light transmittance change.\nFurthermore, since the cyanine compound has a low light resistance in general, both the near-infrared absorptivity and the visible light transmittance of a near-infrared shield using the cyanine alone will change.\nIn addition, a near-infrared shield used in a display front plate is preferred to have excellent near-infrared shielding property and visible light transmittance, and be capable of enduring long-term storage under a condition of high temperature, high humidity and light irradiation."}
{"text": "Conventionally, an axle driving apparatus consists of a housing for an HST, axles and a power transmitting device for interconnecting the HST and axles. On the center section of the HST is disposed a hydraulic pump, provided with a vertical input shaft, and a hydraulic motor, provided with a horizontal output shaft. A plurality of pistons are disposed in the hydraulic pump cylinder block. The heads of the pistons abut against a movable swash plate. Changing the angle of the movable swash plate changes the pump capacity so as to increase or decrease the number of rotations of the hydraulic motor. The movable swash plate is slanted, thereby enabling the speed of the HST to be changed by rotatably operating trunnions supported in the housing. Each trunnion is disposed on a longitudinally slanted axis of the swash plate, as disclosed in U.S. Pat. No. 5,456,068, for example.\nA speed change controller, such as a pedal or a lever, which is provided on the vehicle can be operated normally longitudinally thereof so that its motion can be transmitted to a control arm of the axle driving apparatus through a link mechanism, such as a rod, disposed longitudinally of the vehicle. Hence, it is preferable that the control arm swing longitudinally around the lateral axis. One conventional construction is provided with a vertical operating shaft, independent of the trunnions, where both trunnions and the vertical operating shaft interlock with each other. The control arm is provided at one end of the operating shaft so that the control arm swings longitudinally around the vertical axis, and the other end is constructed so that the trunnion projects at the axial end thereof from the front wall of the housing. A control arm is provided at the axial end so that the control arm swings laterally around the longitudinal axis. A complex linkage mechanism, with respect to the vertical operating shaft and trunnions, is required in the first construction described above, thereby increasing the number of parts and assembly time, making the axle driving apparatus too expensive to produce. The second construction described above requires a separate link mechanism for converting the longitudinal motion into a lateral motion, thereby requiring space to provide two link mechanisms in the vehicle, making it difficult to apply the apparatus to a vehicle of small size and increasing the number of parts required.\nU.S. Pat. Nos. 5,440,951 and 5,515,747 disclose that when the HST and the mechanism for transmitting power to the axles from the HST are housed in the same housing, the housing can be filled with oil to be used as both operating oil for the HST and lubricating oil for the transmitting mechanism. In this case, a foreign object, such as iron powder, created by the rubbing of the transmitting mechanism may flow toward the HST. The iron powder or other foreign object is removed by an oil filter so as not to enter into the HST closed fluid circuit. However, the iron powder or the like may encroach on the piston and swash plate and thereby adversely affect them. The housing is integrated in part with the oil reservoir so as to enable the oil volume in the housing to be adjusted when expanded due to a rise in temperature. However, the greater the quantity of oil, the larger the increase in volume. Thus, the housing must be made larger and the reservoir therefore becomes larger so that the housing itself has to be large in size.\nU.S. Pat. No. 5,094,077 discloses that in order to prevent the speed change controller equipped on the vehicle from being hastily operated by an operator, a shock absorber is provided on the control arm. The shock absorber should be disposed above the upper wall of the housing because the control arm is configured to vertically and longitudinally swing around the axis on the upper wall of the housing. Therefore, space for disposing the shock absorber without interference with an input pulley or an enlarged portion of the upper wall of the housing is required.\nFurther, where a differential gear is provided between the left and right axles, when one axle is idling, a driving force cannot be transmitted to the other axle. Hence, it is desired to provide a differential locking device on the axle driving apparatus for integrating the differential locking device with the HST and the axles.\nAdditionally, conventionally there is a well-known IHT, which comprises a housing containing an HST, a pair of axles and a differential unit. A problem arises in the IHT having the differential unit interposed between left and right axles. For example, a vehicle equipped with the IHT, when one of left and right drive wheels is mired in mud or a ditch, cannot escape because the mired wheel idles therein so as to hinder the other wheel from receiving power."}
{"text": "The present invention relates to catalyst components for the polymerization of olefins, in particular propylene, comprising a Mg dihalide based support on which are deposited a Ti compound having at least one Ti-halogen bond and at least two electron donor compounds selected from specific classes. The present invention further relates to the catalysts obtained from said components and to their use in processes for the polymerization of olefins. The catalysts of the present invention are able to give, with high yields, polymers characterized by high xylene insolubility, a broad range of isotacticity and are further characterized by a good balance between hydrogen response and isotacticity.\nCatalyst components for the stereospecific polymerization of olefins are widely known in the art. Basically two types of catalyst systems are used in the normal processes for the (co)polymerization of olefins. The first one, in its broadest definition, comprises TiCl3 based catalysts components, obtained for example by reduction of TiCl4 with Al-alkyls, used in combination with Al-compounds such as diethylaluminum chloride (DEAC). Despite the good properties of the polymers in terms of isotacticity said catalysts are characterized by a very low activity which causes the presence of large amounts of catalytic residues in the polymers. As a consequence, a further step of deashing is necessary to obtain a polymer having a content of catalytic residue that makes it acceptable for wide use.\nThe second type of catalyst system comprises a solid catalyst component, constituted by a magnesium dihalide on which are supported a titanium compound and an internal electron donor compound, used in combination with an Al-alkyl compound. Conventionally however, when a higher crystallinity of the polymer is required, also an external donor (for example an alkoxysilane) is needed in order to obtain higher isotacticity. One of the preferred classes of internal donors is constituted by the esters of phthalic acid, diisobutylphthalate being the most used. This catalyst system is capable to give very good performances in terms of activity, isotacticity and xylene insolubility provided that an external electron donor compound is used. In its absence, low yields, low xylene insolubility and poor isotacticity are obtained. On the other hand, when the external donor is used, high xylene insolubility is obtained only together with a high isotacticity. This is not desirable in certain applications, such as production of bi-oriented polypropylene films (BOPP), where polypropylenes are required to have a lower flexural modulus (obtainable by lowering crystallinity of the polymer) while at the same time retaining a high xylene insolubility. As a consequence, it would be desirable to have a catalyst component with still improved characteristics, particularly in terms of activity and isotacticity, as well a catalyst component capable to give polymers coupling high xylene insolubility with a slight lower crystallinity suitable for making the polymers usable in the BOPP sector. Some improvements are obtained when, in the above mentioned catalyst system, the phthalates are substituted by the electron donor compounds disclosed for example in U.S. Pat. No. 4,971,937. In this case, the catalyst components obtained are capable to give better results when used in the absence of an external donor. In particular, the stereoregularity becomes acceptable, while however the xylene insolubility is still to be improved. Also in this case, when the catalyst component is used together with an external donor, high xylene insolubility isnobtaied only together with a high isotacticity.\nIt is therefore felt the need of a versatile catalyst component which, for high values of xylene insolubility, is capable to give polymers with a broader range of isotacticity. Moreover, it/would be also advantageous to have a catalyst component with still improved features in terms of activity and isotacticity.\nIt has now unexpectedly been found a catalyst component having the above advantages which comprises Mg, Ti, halogen and two electron donor compounds selected from specific classes.\nIt is therefore an object of the present invention a catalyst component for the polymerization of olefins CH2xe2x95x90CHR, in which R is hydrogen or a hydrocarbyl radical with 1-12 carbon atoms, comprising Mg, Ti, halogen and at least two electron donor compounds, said catalyst component being characterized by the fact that at least one of the electron donor compounds is selected from ethers containing two or more ether groups which are/further characterized by the formation of complexes with anhydrous magnesium dichloride in an amount less than 60 mmoles per 100 g of MgCl2 and by the failure of entering into substitution reactions with TiCl4 or by reacting in that way for less than 50% by moles, and at least another electron donor compound is selected from esters of mono or polycarboxylic acids.\nThe conditions under which, the reactivity toward titanium tetrachloride and the complexing activity of the di or polyethers are tested, are reported below.\nVery surprisingly it has been found that the performances of the above-disclosed catalysts are not merely intermediate between those of the catalyst components containing the single donors. While we do not intend being bound to any theoretical interpretation, it can be said that a synergic interaction between the elements of the catalyst component, and maybe in particular between the above mentioned donors, is the basis for explaining the unexpected properties of the catalyst component of the invention.\nAmong the di or polyethers mentioned above, particularly preferred are the compounds belonging to the class of the 1,3-diethers. In particular, preferred 1,3-diethers are those of formula (I) \nwhere RI and RII are the same or different and are hydrogen or linear or branched C1-C18 hydrocarbon groups which can also form one or more cyclic structures; RIII groups, equal or different from each other, are hydrogen or C1-C18 hydrocarbon groups; RIV groups equal or different from each other, have the same meaning of RIII except that they cannot be hydrogen; each of RI to RIV groups can contain heteroatoms selected from halogens, N, O, S and Si.\nPreferably, RIV is a 1-6 carbon atom alkyl radical and more particularly a methyl while the RIII radicals are preferably hydrogen. Moreover, when RI is methyl, ethyl, propyl, or isopropyl, RII can be ethyl, propyl, isopropyl, butyl, isobutyl, tert-butyl, isopentyl, 2-ethyihexyl, cyclopentyl, cyclohexyl, methylcyclohexyl, phenyl or benzyl; when RI is hydrogen, RII can be ethyl, butyl, sec-butyl, tert-butyl, 2-ethylhexyl, cyclohexylethyl, diphenylmethyl, p-chlorophenyl, 1-naphthyl, 1-decahydronaphthyl; RI and RII can also be the same and can be ethyl, propyl, isopropyl, butyl, isobutyl, tert-butyl, neopentyl, phenyl, benzyl, cyclohexyl, cyclopentyl.\nSpecific examples of ethers that can be advantageously used include: 2-(2-ethylhexyl)1,3-dimethoxypropane, 2-isopropyl-1,3-dimethoxypropane, 2-butyl-1,3-dimethoxypropane, 2-sec-butyl-1,3-dimethoxypropane, 2-cyclohexyl-1,3-dimethoxypropane, 2-phenyl-1,3-dimethoxypropane, 2-tert-butyl-1,3-dimethoxypropane, 2-cumyl-1,3-dimethoxypropane, 2-(2-phenylethyl)-1,3-dimethoxypropane, 2-(2-cyclohexylethyl)-1,3-dimethoxypropane, 2-(p-chlorophenyl)-1,3-dimethoxypropane, 2-(diphenylmethyl)-1,3-dimethoxypropane, 2(1-naphthyl)-1,3-dimethoxypropane, 2(p-fluorophenyl)-1,3-dimethoxypropane, 2(1-decahydronaphthyl)-1,3-dimethoxypropane, 2(p-tert-butylphenyl)-1,3-dimethoxypropane, 2,2-dicyclohexyl-1,3-dimethoxypropane, 2,2-diethyl-1,3-dimethoxypropane, 2,2-dipropyl-1,3-dimethoxypropane, 2,2-dibutyl-1,3-dimethoxypropane, 2,2-diethyl-1,3-diethoxypropane, 2,2-dicyclopentyl-1,3-dimethoxypropane, 2,2-dipropyl-1,3-diethoxypropane, 2,2-dibutyl-1,3-diethoxypropane, 2-methyl-2-ethyl-1,3-dimethoxypropane, 2-methyl-2-propyl-1,3-dimethoxypropane, 2-methyl-2-benzyl-1,3-dimethoxypropane, 2-methyl-2-phenyl-1,3-dimethoxypropane, 2-methyl-2-cyclohexyl-1,3-dimethoxypropane, 2-methyl-2-methylcyclohexyl-1,3-dimethoxypropane, 2,2-bis(p-chlorophenyl)-1,3-dimethoxypropane, 2,2-bis(2-phenylethyl)-1,3-dimethoxypropane, 2,2-bis(2-cyclohexylethyl)-1,3-dimethoxypropane, 2-methyl-2-isobutyl-1,3-dimethoxypropane, 2-methyl-2-(2-ethylhexyl)-1,3-dimethoxypropane, 2,2-bis(2-ethylhexyl)-1,3-dimethoxypropane,2,2-bis(p-methylphenyl)-1,3-dimethoxypropane, 2-methyl-2-isopropyl-1,3-dimethoxypropane, 2,2-diisobutyl-1,3-dimethoxypropane, 2,2-diphenyl-1,3-dimethoxypropane, 2,2-dibenzyl-1,3-di methoxypropane, 2-isopropyl-2-cyclopentyl-1,3-dimethoxypropane, 2,2-bis(cyclohexylmethyl)-1,3-dimethoxypropane, 2,2-diisobutyl-1,3-diethoxypropane, 2,2-diisobutyl-1,3-dibutoxypropane, 2-isobutyl-2-isopropyl-1,3-dimetoxypropane, 2,2-di-sec-butyl-1,3-dimetoxypropane, 2,2-di-tert-butyl-1,3-dimethoxypropane, 2,2-dineopentyl-1,3-dimethoxypropane, 2-iso-propyl-2-isopentyl-1,3-dimethoxypropane, 2-phenyl-2-benzyl-1,3-dimetoxypropane, 2-cyclohexyl-2-cyclohlexylmethyl-1,3-dimethoxypropane.\nFurthermore, particularly preferred are the 1,3-diethers of formula (II) \nwhere the radicals RIV have the same meaning explained above and the radicals RIII and RV radicals, equal or different to each other, are selected from the group consisting of hydrogen; halogens, preferably Cl and F; C1-C20 alkyl radicals, linear or branched; C3-C20 cycloalkyl, C6-C20 aryl, C7-C20 alkaryl and C7-C20 aralkyl radicals and two or more of the RV radicals can be bonded to each other to form condensed cyclic structures, saturated or unsaturated, optionally substituted with RVI radicals selected from the group consisting of halogens, preferably Cl and F; C1-C20 alkyl radicals, linear or branched; C1-C20 cycloalkyl, C6-C20 aryl, C7-C20 alkaryl and C7-C20 aralkyl radicals; said radicals RV and RVI optionally containing one or more heteroatoms as substitutes for carbon or hydrogen atoms, or both. Preferably, in the 1,3-diethers of formulae (I) and (II) all the RIII radicals are hydrogen, and all the RIV radicals are methyl. Moreover, are particularly preferred the 1,3-diethers of formula (II) in which two or more of the RV radicals are bonded to each other to form one or more condensed cyclic structures, preferably benzenic, optionally substituted by RVI radicals. Specially preferred are the compounds of formula (III): \nwhere the RVI radicals equal or different are hydrogen; halogens, preferably Cl and F; C1-C20 alkyl radicals, linear or branched; C3-C20 cycloalkyl, C6-C20 aryl, C7-C20 alkylaryl and C7-C20 aralkyl radicals, optionally containing one or more heteroatoms selected from the group consisting of N, O, S, P, Si and halogens, in particular Cl and F, as substitutes for carbon or hydrogen atoms, or both; the radicals RIII and RIV are as defined above for formula (II).\nSpecific examples of compounds comprised in formulae (II) and (III) are:\n1,1-bis(methoxymethyl)-cyclopentadiene;\n1,1-bis(methoxymethyl)-2,3,4,5-tetramethylcyclopentadiene;\n1,1-bis(methoxymethyl)-2,3,4,5-tetraphenylcyclopentadiene;\n1,1-bis(methoxymethyl)-2,3,4,5-tetrafluorocyclopentadiene;\n1,1-bis(methoxymethyl)-3,4-dicyclopentylcyclopentadiene;\n1,1-bis(methoxymethyl)indene; 1,1-bis(methoxymethyl)-2,3-dimethylindene;\n1,1-bis(imethoxymethyl)-4,5,6,7-tetrahydroindene;\n1,1-bis(methoxymethyl)-2,3,6,7-tetrafluoroindene;\n1,1-bis(methoxymethyl)-4,7-dimethylindene;\n1,1-bis(methoxymethyl)-3,6-dimethylindene;\n1,1-bis(methoxymethyl)-4-phenylindene;\n1,1-bis(methoxymethyl)-4-phenyl-2-methylindene;\n1,1-bis(methoxymethyl)-4-cyclohexylindene;\n1,1-bis(methoxymethyl)-7-(3,3,3-trifluoropropyl)indene;\n1,1-bis(methoxymethyl)-7-trimethyisilylindene;\n1,1-bis(methoxymethyl)-7-trifluoromethylindene;\n1,1-bis(methoxymethyl)-4,7-dimethyl-4,5,6,7-tetrahydroindene;\n1,1-bis(methoxymethyl)-7-methylindene;\n1,1-bis(methoxymethyl)-7-cyclopenthylindene;\n1,1-bis(methoxymethyl)-7-isopropylindene;\n1,1-bis(methoxymethyl)-7-cyclohexylindene;\n1,1-bis(methoxymethyl)-7-tert-butylindene;\n1,1-bis(methoxymethyl)-7-tert-butyl-2-methylindene;\n1,1-bis(methoxymethyl)-7-phenylindene;\n1,1-bis(methoxymethyl)-2-phenylindene;\n1,1-bis(methoxymethyl)-1H-benz[e]indene;\n1,1-bis(methoxymethyl)-1H-2-methylbenz[e]indene;\n9,9-bis(methoxymethyl)fluorene;\n9,9-bis(methoxymethyl)-2,3,6,7-tetramethylfluorene;\n9,9-bis(methoxymethyl)-2,3,4,5,6,7-hexafluorofluorene;\n9,9-bis(methoxymethyl)-2,3-benzofluorene;\n9,9-bis methoxymethyl)-2,3,6,7-dibenzofluorene;\n9,9-bis(methoxymethyl)-2,7-diisopropylfluorene;\n9,9-bis(methoxymethyl)-1,8-dichlorofluorene;\n9,9-bis(methoxymethyl)-2,7-dicyclopentylfluorene;\n9,9-bis(methoxymethyl)-1,8-difluorofluorene;\n9,9-bis(methoxymethyl)-1,2,3,4-tetrahydrofluorene;\n9,9-bis(methoxymethyl)-1,2,3,4,5,6,7,8-octahydrofluorene;\n9,9-bis(methoxymethyl)-4-tert-butylfluorene.\nAs explained above, the other electron donor compound which must be present in the catalyst component of the invention has to be selected from the esters of mono or polycarboxylic acids. Said acids can be both aliphatic and aromatic acids.\nAmong esters of aliphatic acids, particularly preferred are the esters of bicarboxylic acids in particular esters of malonic acids. Particularly preferred are the esters of malonic acids of formula (IV): \nwhere R1 is H or a C1-C20 linear or branched alkyl, alkenyl, cycloalkyl, aryl, arylalkyl or alkylaryl group, R2 is a C1-C20 linear or branched alkyl, alkenyl, cycloalkyl, aryl, arylalkyl or alkylaryl group, R3 and R4, equal to, or different from, each other, are C1-C20 linear or branched alkyl groups or C3-C20 cycloalkyl groups.\nPreferably, R3 and R4 are primary, linear or branched C1-C20 alkyl groups, more preferably they are primary branched C4-C20 alkyl groups such as isobutyl or neopentyl groups. R2 is preferably, in particular when R1 is H, a linear or branched C3-C20 alkyl, cycloalkyl, or arylalkyl group; more preferably R2 is a C3-C20 secondary alkyl, cycloalkyl, or arylalkyl group.\nSpecific examples of preferred monosubstituted malonate compounds are: dineopentyl 2-isopropylmalonate, diisobutyl 2-isopropylmalonate, di-n-butyl 2-isopropylmalonate, diethyl 2-dodecylmalonate, diethyl 2-t-butylmalonate, diethyl 2-(2-pentyl)malonate, diethyl 2-cyclohexylmalonate, dineopentyl 2-t-butylmalonate, dineopentyl 2-isobutylmalonate, diethyl 2-cyclohexylmethylmalonate, dimethyl 2-cyclohexylmethylmalonate.\nSpecific examples of preferred disubstituted malonates compounds are: diethyl 2,2-dibenzylmalonate, diethyl 2-isobutyl-2-cyclohexylmalonate, dimethyl 2-n-butyl-2-isobutylmalonate, diethyl 2-n-butyl-2-isobutylmalonate, diethyl 2-isopropyl-2-n-butylmalonate, diethyl 2-methyl-2-isopropylmalonate, diethyl 2-isopropyl-2-isobutylmalonate, diethyl 2-methyl-2-isobutylmalonate, diethyl 2-isobutyl-2-benzylmalonate.\nPreferred esters of aromatic carboxylic acids are selected from C1-C20 alkyl or aryl esters of benzoic and phthalic acids, possibly substituted. The alkyl esters of the said acids being preferred. Particularly preferred are the C1-C6 linear or branched alkyl esters. Specific examples are ethylbenzoate, n-butylbenzoate, p-methoxy ethylbenzoate, p-ethoxy ethylbenzoate, isobutylbenzoate, ethyl p-toluate, diethyl phthalate, di-n-propyl phthalate, di-n-butyl phthalate di-n-pentyl phthalate, di-i-pentyl phthalate, bis(2-ethylhexyl) phthalate, ethyl-isobutyl phthalate, ethyl-n-butyl phthalate, di-n-hexyl phthalate, di-isobutylphthalate.\nAs explained above, the catalyst components of the invention comprise, in addition to the above electron donors, Ti, Mg and halogen. In particular, the catalyst components comprise a titanium compound, having at least a Ti-halogen bond and the above mentioned electron donor compounds supported on a Mg halide. The magnesium halide is preferably MgCl2 in active form which is widely known from the patent literature as a support for Ziegler-Natta catalysts. Patents U.S. Pat. No. 4,298,718 and U.S. Pat. No. 4,495,338 were the first to describe the use of these compounds in Ziegler-Natta catalysis. It is known from these patents that the magnesium dihalides in active form used as support or co-support in components of catalysts for the polymerization of olefins are characterized by X-ray spectra in which the most intense diffraction line that appears in the spectrum of the non-active halide is diminished in intensity and is replaced by a halo whose maximum intensity is displaced towards lower angles relative to that of the more intense line.\nThe preferred titanium compounds used in the catalyst component of the present invention are TiCl4 and TiCl3; furthermore, also Ti-haloalcoholates of formula Ti(OR)nxe2x88x92yXy can be used, where n is the valence of titanium, y is a number between 1 and nxe2x88x921 X is halogen and R is a hydrocarbon radical having from 1 to 10 carbon atoms.\nThe preparation of the solid catalyst component can be carried out according to several methods. According to one of these methods, the magnesium dichloride in an anhydrous state, the titanium compound and the electron donor compounds are milled together under conditions in which activation of the magnesium dichloride occurs. The so obtained product can be treated one or more times with an excess of TiCl4 at a temperature between 80 and 135xc2x0 C. This treatment is followed by washings with hydrocarbon solvents until chloride ions disappeared. According to a further method, the product obtained by co-milling the magnesium chloride in an anhydrous state, the titanium compound and the electron donor compounds are treated with halogenated hydrocarbons such as 1,2-dichloroethane, chlorobenzene, dichloromethane etc. The treatment is carried out for a time between 1 and 4 hours and at temperature of from 40xc2x0 C. to the boiling point of the halogenated hydrocarbon. The product obtained is then generally washed with inert hydrocarbon solvents such as hexane.\nAccording to another method, magnesium dichloride is preactivated according to well known methods and then treated with an excess of TiCl4 at a temperature of about 80 to 135xc2x0 C. in the presence of the electron donor compounds. The treatment with TiCl4 is repeated and the solid is washed with hexane in order to eliminate any non-reacted TiCl4. A further method comprises the reaction between magnesium alcoholates or chloroalcoholates (in particular chloroalcoholates prepared according to U.S. Pat. No. 4,220,554) and an excess of TiCl4 in the presence of the electron donor compounds at a temperature of about 80 to 120xc2x0 C.\nAccording to a preferred method, the solid catalyst component can be prepared by reacting a titanium compound of formula Ti(OR)nxe2x88x92yXy, where n is the valence of titanium and y is a number between 1 and n, preferably TiCl4, with a magnesium chloride deriving from an adduct of formula MgCl2pROH, where p is a number between 0.1 and 6, preferably from 2 to 3.5, and R is a hydrocarbon radical having 1-18 carbon atoms. The adduct can be suitably prepared in spherical form by mixing alcohol and magnesium chloride in the presence of an inert hydrocarbon immiscible with the adduct, operating under stirring conditions at the melting temperature of the adduct (100-130xc2x0 C.). Then, the emulsion is quickly quenched, thereby causing the solidification of the adduct in form of spherical particles. Examples of spherical adducts prepared according to this procedure are described in U.S. Pat. No. 4,399,054 and U.S. Pat. No. 4,469,648. The so obtained adduct can be directly reacted with Ti compound or it can be previously subjected to thermal controlled dealcoholation (80-130xc2x0 C.) so as to obtain an adduct in which the number of moles of alcoholis generally lower than 3 preferably between 0.1 and 2.5. The reaction with the Ti compound can be carried out by suspending the adduct (dealcoholated or as such) in cold TiCl4 (generally 0xc2x0 C.); the mixture is heated up to 80-130xc2x0 C. and kept at this temperature for 0.5-2 hours. The treatment with TiCl4 can be carried out one or more times. The electron donor compounds can be added during the treatment with TiCl4. They can be added together in the same treatment with TiCl4 or separately in two or more treatments. The preparation of catalyst components in spherical form are described for example in European Patent Applications EP-A-395083, EP-A-553805, EP-A-553806, EPA601525 and WO98/44001.\nThe solid catalyst components obtained according to the above method show a surface area (by B.E.T. method) generally between 20 and 500 m2/g and preferably between 50 and 400 m2/g and a total porosity (by B.E.T. method) higher than 0.2 cm3/g preferably between 0.2 and 0.6 cm3/g. The porosity (Hg method) due to pores with radius up to 10.000 xc3x85 generally ranges from 0.3 to 1.5 cm3/g, preferably from 0.45 to 1 cm3/g.\nA further method to prepare the solid catalyst component of the invention comprises halogenating magnesium dihydrocarbyloxide compounds, such as magnesium dialkoxide or diaryloxide, with solution of TiCl4 in aromatic hydrocarbon (such as toluene, xylene etc.) at temperatures between 80 and 130xc2x0 C. The treatment with TiCl4 in aromatic hydrocarbon solution can be repeated one or more times, and the electron donor compounds are added during one or more of these treatments.\nIn any of these preparation methods the desired electron donor compounds and in particular those selected from esters of carboxylic acids, can be added as such or, in an alternative way, it can be obtained in situ by using an appropriate precursor capable to be transformed in the desired electron donor compound by means, for example, of known chemical reactions such as esterification, transesterification, etc.\nRegardless to the preparation method used, the final amount of the two or more electron donor compounds is such that the molar ratio with respect to the MgCl2 is from 0.01 to 1, preferably from 0.05 to 0.5, while the molar ratio between the di or polyether donor and the esteddonor is comprised in the range of from 50 to 0.02 preferably from 30 to 0.1 and more preferably from 20 to 0.2.\nThe solid catalyst components according to the present invention are converted into catalysts for the polymerization of olefins by reacting them with organoaluminum compounds according to known methods.\nIn particular, it is an object of the present invention a catalyst for the polymerization of olefins CH2xe2x95x90CHR, in which R is hydrogen or a hydrocarbyl radical with 1-12 carbon atoms, comprising the product of the reaction between:\n(i) the solid catalyst component as disclosed above and\n(ii) an alkylaluminum compound.\nThe alkyl-Al compound (ii) is preferably chosen among the trialkyl aluminum compounds such as for example triethylaluminum, triisobutylaluminum, tri-n-butylaluminum, tri-n-hexylaluminum, tri-n-octylaluminum. It is also possible to use alkylaluminum halides, alkylaluminum hydrides or alkylaluminum sesquichlorides, such as AlEt2Cl and Al2Et3Cl3, possibly in mixture with the above cited trialkylaluminums.\nAs explained above, the catalyst component of the invention when used in the polymerization of propylene in the absence of external donors are able to give polymers with a controlled wide range of isotacticity (expressed in term of percentage of mmmm pentads) while maintaining high xylene insolubility levels.\nFurthermore, the catalyst components of the invention can also be used in combination with an external donor (iii) thereby obtaining very high values of both xylene insolubility and isotacticity. In particular, said values, individually or as a balance, are higher than the values obtainable with the catalyst containing the single donors.\nSuitable external electron-donor compounds include silicon compounds, ethers, esters, amines, heterocyclic compounds and particularly 2,2,6,6-tetramethyl piperidine, ketones and the 1,3-diethers of the general formula (V): \nwherein RI, RII, RIII, RIV, RV and RVI equal or different to each other, are hydrogen or hydrocarbon radicals having from 1 to 18 carbon atoms, and RVII and RVIII, equal or different from each other, have the same meaning of RI-RVI except that they cannot be hydrogen; one or more of the RI-RVIII groups can be linked to form a cycle. Particularly preferred are the 1,3-diethers in which RVII and RVIII are selected from C1-C4 alkyl radicals.\nAnother class of preferred external donor compounds is that of silicon compounds of formula Ra5Rb6Si(OR7)c, where a and b are integers from 0 to 2, c is an integer from 1 to 3 and the sum (a+b+c) is 4; R5, R6, and R7, are alkyl, cycloalkyl or aryl radicals with 1-18 carbon atoms optionally containing heteroatoms. Particularly preferred are the silicon compounds in Which a is 1, b is 1, c is 2, at least one of R5 and R6 is selected from branched alkyl, cycloalkyl or aryl groups with 3-10 carbon atoms optionally containing heteroatoms and R7 is a C1-C10 alkyl group, in particular methyl. Examples of such preferred silicon compounds are methylcyclohexyldimethoxysilane, diphenyldimethoxysilane, methyl-t-butyldimethoxysilane, dicyclopentyldimethoxysilane, (2-ethylpiperidinyl)t-butyldimethoxysilane, (2-ethylpiperidinyl)thexyldimethoxysilane, (3,3,3-trifluoro-n-propyl)(2-ethylpiperidinyl)dimethoxysilane, methyl(3,3,3-trifluoro-n-propyl)dimethoxysilane. Moreover, are also preferred the silicon compounds in which a is 0, c is 3, R6 is abranched alkyl or cycloalkyl group, optionally containing heteroatoms, and R7 is methyl. Examples of such preferred silicon compounds are cyclohexyltrimethoxysilane, t-butyltrimethoxysilane and thexyltrimethoxysilane.\nThe electron donor compound (iii) is used in such an amount to give a molar ratio between the organoaluminum compound and said electron donor compound (iii) of from 0.1 to 500, preferably from 1 to 300 and more preferably from 3 to 100.\nTherefore, it constitutes a further object of the present invention a process for the (co)polymerization of olefins CH2xe2x95x90CHR, in which R is hydrogen or a hydrocarbyl radical with 1-12 carbon atoms, carried out in the presence of a catalyst comprising the product of the reaction between:\n(i) a solid catalyst component comprising a titanium compound, having at least a Ti-halogen bond, and at least two electron donor compounds supported on a Mg halide said catalyst being characterized by the fact that at least one of the electron donors compounds is selected from ethers containing two or more ether groups which are further characterized by the formation of complexes with anhydrous magnesium dichloride in an amount less than 60 mmoles per 100 g of MgCl2 and by the failure of entering into substitution reactions with TiCl4 or by reacting in that way for less than 50% by moles, and at least another electron donor compound is selected from esters of mono or polycarboxylic acids;\n(ii) an alkylaluminum compound and,\n(iii) optionally an electron-donor compound (external donor).\nThe polymerization process can be carried out according to known techniques for example slurry polymerization using as diluent an inert hydrocarbon solvent, or bulk polymerization using the liquid monomer (for example propylene) as a reaction medium. Moreover, it is possible to carry out the polymerization process in gas-phase operating in one or more fluidized or mechanically agitated bed reactors.\nThe polymerization is generally carried out at temperature of from 20 to 120xc2x0 C., preferably of from 40 to 80xc2x0 C. When the polymerization is carried out in gas-phase the operating pressure is generally between 0.5 and 5 MPa, preferably between 1 and 4 MPa. In the bulk polymerization the operating pressure is generally between 1 and 8 MPa preferably between 1.5 and 5 MPa."}
{"text": "Footwear typically includes a sole configured to be located under a wearer's foot to space the foot away from the ground or floor surface. Soles can be designed to provide a desired level of cushioning. Athletic footwear in particular sometimes utilizes polyurethane foam or other resilient materials in the sole to provide cushioning. Fluid-filled bladders are sometimes included in the sole to provide desired impact force absorption, motion control, and resiliency. The incorporation of additional materials and components adds processing steps to the manufacturing of footwear."}
{"text": "The present invention relates to a novel gene coding for thermostable beta-galactosidase, Bacillus subtilis, having the above gene, a novel enzyme coded by the gene and a process for the production thereof.\nBeta-galactosidase is an enzyme which hydrolyzes lactose to galactose and glucose and is widely utilized in food processing, such as the production of milk with a low lactose content or the production of galactose or glucose from lactose contained in whey obtained in a large amount as a by-product in the production of cheese.\nGenerally, it is desired that enzymes for use in food processing be stable at an elevated temperature from the viewpoint of preventing microbial contamination. This also applies to the aforementioned enzyme.\nFurther, the enzyme is also utilized as a pharmaceutical to treat intolerance to lactose. In this case, good heat resistance is preferred for stability of the preparation as well.\nThe present invention has been made in order to achieve the above various requirements and provides a well-defined process for producing beta-galactosidase of excellent heat stability in a commercially advantageous manner."}
{"text": "1. Field of the Invention\nThe invention relates to a method for modelling steady state as well as transient multiphase flows such as hydrocarbon mixtures circulating in pipeline networks taking into account a set of variables defining fluid properties and flow patterns, as well as dimensions or slope angles of the pipelines.\n2. Description of the Prior Art\nPrior art relating to multi-phase\nFabre, J., et al 1983. Intermittent gas-liquid flow in horizontal or slightly inclined pipes, Int. Conference on the Physical Modelling of Multi-Phase Flow, Coventry, England, pp 233, 254 PA1 Fabre, J., et al 1989. Two fluid/two flow pattern model for transient gas liquid flow in pipes, Int. Conference on Multi-Phase Flow, Nice, France, pp 269, 284, Cranfield, BHRA. PA1 characterizing flow regimes by a parameter .beta. representing the fraction of a flow in a separated state, said parameter .beta. continuously ranging from 0 for dispersed flow regimes and 1 for separated flow regimes; PA1 determining any current flow regime while determining said set of transport equations by comparing current values of a liquid fraction in slugs with respect to a liquid fraction in a dispersed region of the flow, as well as gas slug velocity with respect to a critical velocity; and PA1 imposing while solving said closure relations continuity constraints at the boundaries between said regimes to the respective gas volume fractions and to the slug velocities."}
{"text": "This invention relates to spray shields primarily for use with surgical scissors to protect people from potentially hazardous blood spray when an umbilical cord or any other item is severed.\nFor more than fifty years, labor and delivery clinicians have used the same method for transecting the umbilical cord. The process includes placing two clamps on the umbilical cord several inches apart from each other and cutting the umbilical cord between the two clamps with surgical scissors. The blood in the umbilical cord between the two clamps remains pressurized and has been known to spray as far as eight feet or more hitting walls, ceiling tiles, drapes, clinicians and others. It is well established that there exists more than twenty blood borne pathogens such as AIDS/HIV, Hepatitis C, sexually transmitted diseases and now the Ebola virus that could be present in the blood spray and which could contaminate those present. Cases of clinician exposure to blood borne pathogens such as HIV/AIDS and sexually transmitted diseases through cord blood spray are well documented. Such cases require extensive, costly and physically demanding prophylactic treatment.\nIn spite of numerous technological attempts to address this problem, no commercially or functionally viable devices or solutions are available to clinicians to prevent blood spray when cutting the cord. In the majority of cases, the prior art/devices simply did not function as intended. In one instance, the device attempted to combine functions such as preventing spray while also drawing cord blood into vials. In another example, the device was intended to prevent blood spray while simultaneously cutting and clamping the umbilical cord. Those and other known devices were plagued with operational issues involving one or more of these functions and consequently unable to perform the intended functions to the satisfaction of the market. Devices that required a significant departure from existing procedures and new equipment also met with resistance. As with virtually every decision, cost is an important factor and the cost of some devices were considered prohibitive.\nWith no functional solution available, clinicians try to minimize the impact of umbilical cord blood spray in several different ways. Many clinicians put a hand up near the scissors to act as a shield. Unfortunately, this prevents the clinician from using that hand for clinical or care related functions and momentarily distracts the clinician causing them to focus on their own safety rather than the immediate procedure and the patient's care and safety.\nAnother approach is to look away when the cord is cut hoping any spray that might occur will hit them somewhere other than the eyes, nose or mouth. Clearly this is undesirable as it is best to have clinicians looking at the task at hand.\nAlternatively, some caregivers have even held a towel over the area in an attempt to block the spray. None of these approaches are effective or safe for either the clinician or the patient. Furthermore, they do not comply with the federal requirements as set forth in 29 CFR 1910 generally requiring that employers reduce the risk of blood borne pathogen exposure to employees.\nOne attempt to address this problem is disclosed as a disposable shield in U.S. Pat. No. 5,542,435 to Kelly, et al. This device was intended to protect only the scissor user. It does not protect any other personnel in the area due to its shape, orientation, location and distance from the source of the spray.\nWhile protecting the clinician cutting the cord, typically a doctor or midwife, is a worthwhile endeavor, a 2012 survey of labor and delivery nurses (i.e. clinicians not cutting the umbilical cord and positioned in physical locations throughout the delivery area) revealed that 95% of those nurses had experienced cord blood spray and 21% had been sprayed within the prior year. Additionally, tests designed to replicate spray from an umbilical cord were recently conducted and the spray pattern analysis indicated that as much as 50% of the spray occurs in a direction other than towards the user. Blood spray is a serious problem to all in the delivery room and is not limited to the clinician cutting the umbilical cord.\nA similar potential solution is disclosed in U.S. Pat. No. D399,971 to Scherer which approaches the problem in a similar manner as does Kelly, et al., but with a different shield design. Nevertheless, Scherer fails to solve the above-noted shortcomings of Kelly, et al.\nAs previously mentioned, there is currently no functionally or commercially viable solution in the marketplace for the problem of blood spray resulting from the transection of the umbilical cord or other cutting procedures."}
{"text": "The invention relates to a transmitting and/or receiving arrangement for portable appliances, consisting of a shielding housing of metal, containing the radio-frequency section, and an antenna, the antenna consisting of two or more antenna resonators which are parasitically coupled to one another and are essentially identically oriented in the longitudinal direction, having in each case one free resonator end and in each case one end angled via a bending edge and conductively connected to the shielding housing.\nIn EP-A 177 362, a wide-band antenna for portable radio appliances is described. It consists of two angled resonators of different resonant frequency. The two resonators are fed by a common line via a branch of the type of an \"inverted-F antenna\". The antenna resonators act independently of one another and do not form a unit. This is why the efficiency is not particularly high. The distance between the parallel legs is constant.\nIn U.S. Pat. No. 4,584,585, an antenna with resonators of wire angles is described. One resonator is fed at the end and forms an \"inverted-L antenna\". The second resonator is parasitically coupled. Although this arrangement achieves a good efficiency, the antenna has a fixed impedance and cannot be easily matched. The shape of the active antenna resonator is relatively difficult to bend and, in addition, the baseplate also exhibits a step. The bandwidth is relatively narrow.\nIn GB-A 2 067 842, a microstrip antenna is described which is applied to an insulating ground. Here, too, the distance between the two antenna resonators is constant. The free ends of the resonators are opposite one another. The feedpoint is close to the free end of one resonator.\n\"Inverted-F antennas\" are known from T. Taga and K. Tsunekawa, \"Performance Analysis of a Built-in Planar Inverted F Antenna for 800 MHz Band Portable Radio Units,\" IEEE Trans., Selected Areas in Commun., vol. SAC-5, no.5, pp. 921-929, June (1987). Such antennas are matched by varying the position of the feedpoint."}
{"text": "Currently, a maturely-developed call center technology has been widely applied in various industries. Particularly, with the popularization of a telephone service mode based on the call center technology, not only service level and operation efficiency are greatly improved, but also demands of users for more convenient services are satisfied.\nThe telephone service mode usually includes a manual service and an automatic service. Architecture of a call center system and a basic process of the manual service may be generally described as follows. After a user dials a call center number (for example, China Mobile 10086 or China Post 11185), the call enters the call center system through routing of a communication network; generally, an Automatic Call Distribution (ACD) of the call center is responsible for call access and distribution, and notifies Computer Telephony Integration (CTI) of managing a call access procedure (for example, controlling the call to be accessed by performing functions such as queuing, routing, or connection); the CTI connects, according to a service requested by the call and a current condition of operators, the service with a corresponding operator, and provides the service for the user corresponding to the call through customer service application software.\nIn a basic process of the automatic service, after a call enters the call center system, the CTI and an Interactive Voice Response (IVR) control a service access procedure. The CTI divides different service contents or levels into different queues according to service capabilities and service scopes of agents and operators. The IVR is an important device for processing automatic calls in the call center system, and is capable of providing a voice play for the user, receiving key input selected by the user, and providing different services according to the key information.\nThe service of accessing the call center system may be processed according to many manners, for example, an intelligent service manner, where after the call of the user accesses the call center system, basic information of the user (for example, the level of the user or the account type of the user) is identified according to the number, and the relatively differentiated service is provided for the user according to the basic information. However, in the existing processing manner, the differentiated service is provided merely according to the static information of the user, which has certain limitations, and more refined services cannot be provided for the user."}
{"text": "This section provides background information and introduces information related to various aspects of the disclosure that are described and/or claimed below. These background statements are not admissions of prior art.\nThe usefulness of gallium nitride (GaN), and its ternary and quaternary compounds incorporating aluminum and indium (AlGaN, InGaN, AlInGaN), has been well established, for example, in the fabrication of visible and ultraviolet optoelectronic devices and high-power electronic devices. These devices are typically grown epitaxially using growth techniques including molecular beam epitaxy (MBE), CVD, metalorganic chemical vapor deposition (MOCVD), and hydride vapor phase epitaxy (HYPE).\nAs shown in FIG. 1, GaN and its alloys are most stable in the hexagonal wurtzite crystal structure, in which the structure is described by two (or three) equivalent basal plane axes that are rotated 120° with respect to each other (the a-axes), all of which are perpendicular to a unique c-axis. FIG. 1 illustrates an example c-plane 2, m-plane 4, and a-plane 6. Group III and nitrogen atoms occupy alternating c-planes along the crystal's c-axis. The symmetry elements included in the wurtzite structure dictate that III-nitrides possess a bulk spontaneous polarization along this c-axis, and the wurtzite structure exhibits piezoelectric polarization.\nCurrent nitride technology for electronic and optoelectronic devices employs nitride films grown along the polar c-direction. However, related art c-plane quantum well structures in III-nitride (III-N) based optoelectronic and electronic devices suffer from the undesirable quantum-confined Stark effect (QCSE), due to the existence of strong piezoelectric and spontaneous polarizations. The strong built-in electric fields along the c-direction can significantly degrade the usefulness of these III-N materials.\nOne approach to eliminating the spontaneous and piezoelectric polarization effects in GaN optoelectronic devices is to grow the devices on non-polar planes, the m-planes and a-planes of the crystal. Such planes contain equal numbers of Ga and N atoms and are charge-neutral. Furthermore, subsequent non-polar layers are equivalent to one another so the bulk crystal will not be polarized along the growth direction. However, growth of GaN semiconductor wafers with a non-polar surface remains difficult. Accordingly, there exists a need to increase the efficiency and improve operating characteristics for III-nitride based optoelectronic and electronic devices, for example LEDs."}
{"text": "The detection, analysis, transcription, and amplification of nucleic acids are the most important procedures in modern molecular biology. The application of such procedures for RNA analysis is especially important in the investigation of gene expression, diagnosis of infectious agents or genetic diseases, the generation of cDNA, and analysis of retroviruses, to name but a few applications. The reverse transcription of RNA, followed by polymerase chain reaction amplification, commonly referred to as RT-PCR, has become widely used for the detection and quantification of RNA.\nThe RT-PCR procedure involves two separate molecular syntheses: (i) the synthesis of cDNA from an RNA template; and (ii) the replication of the newly synthesized cDNA through PCR amplification. RT-PCR can be performed under three general protocols: (1) uncoupled RT-PCR, also referred to as two-step RT-PCR; (2) single enzyme coupled RT-PCR (coupled RT-PCR is also referred to as one-step RT-PCR or continuous RT-PCR), in which a single polymerase is used for both the cDNA generation from RNA as well as subsequent DNA amplification; and (3) two (or more) enzyme coupled RT-PCR, in which at least two separate polymerases are used for initial cDNA synthesis and subsequent replication and amplification.\nIn uncoupled RT-PCR, reverse transcription is performed as an independent step using buffer and reaction conditions optimal for reverse transcriptase activity. Following cDNA synthesis, an aliquot of the RT reaction product is used as template for PCR amplification with a thermostable DNA Polymerase, such as Taq DNA Polymerase, under conditions optimal for PCR amplification.\nIn coupled RT-PCR, reverse transcription and PCR amplification are combined into a single reaction mixture. Single enzyme RT-PCR utilizes the reverse transcriptase activity of some DNA polymerases, such as Taq DNA Polymerase and Tth DNA polymerase, whereas two-enzyme RT-PCR typically uses a retroviral or bacterial reverse transcriptase (e.g., AMV-RT, MMLV-RT, HIV-RT, EIAV-RT, RAV2-RT, Carboxydothermus hydrogenoformans DNA Polymerase or a mutant, variant or derivative thereof), and a thermostable DNA polymerase (e.g., Taq, Tbr, Tth, Tih, Tfi, Tfl, Pfu, Pwo, Kod, VENT, DEEPVENT, Tma, Tne, Bst, Pho, Sac, Sso, ES4 and others or a mutant, variant or derivative thereof).\nCoupled RT-PCR provides numerous advantages over uncoupled RT-PCR. Coupled RT-PCR requires less handling of the reaction mixture reagents and nucleic acid products than uncoupled RT-PCR (e.g., opening of the reaction tube for component or enzyme addition in between the two reaction steps), and is therefore less labor intensive, reducing the required number of person hours. Coupled RT-PCR also requires less sample, and reduces the risk of contamination (Sellner and Turbett, 1998).\nSingle enzyme coupled RT-PCR, is the simplest RT-PCR procedure to date. This system is expensive to perform, however, due to the amount of DNA polymerase required. In addition, the single enzyme coupled RT-PCR method has been found to be less sensitive than uncoupled RT-PCR (Cusi et al., 1994), and limited to polymerizing nucleic acids of less than one kilobase pair (>1 kb) in length. Two enzyme RT-PCR systems show increased sensitivity over the single enzyme system generally, even when coupled in a single reaction mixture. This effect has been attributed to the higher efficiency of reverse transcriptase in comparison to the reverse transcriptase activity of DNA polymerases (Sellner and Turbett, 1998).\nAlthough the two enzyme coupled RT-PCR system is more sensitive than the uncoupled protocol, reverse transcriptase has been found to interfere with DNA polymerase during the replication of the cDNA, thus reducing the sensitivity and efficiency of this technique (Sellner et al., 1992; Aatsinki et al., 1994; Mallet et al., (1995)). A variety of solutions to overcome the inhibitory activity of reverse transcriptase on DNA polymerase have been tried, including: increasing the amount of template RNA, increasing the ratio of DNA polymerase to reverse transcriptase, adding modifier reagents that can reduce the inhibitory effect of reverse transcriptase on DNA polymerase (e.g., non-homologous tRNA, T4 gene 32 protein, sulfur or acetate-containing molecules,), and heat-inactivation of the reverse transcriptase before the addition of DNA polymerase.\nAll of these modified RT-PCR methods have significant drawbacks, however. Increasing the amount of template RNA is not possible in cases where only limited amounts of sample are available. Individual optimization of the ratio of reverse transcriptase to DNA polymerase is not practicable for ready-to-use reagent kits for one-step RT-PCR. The net effect of currently proposed modifier reagents to relieve reverse transcriptase inhibition of DNA polymerization is controversial and in dispute: positive effects due to these reagents are highly dependent on RNA template amounts, RNA composition, or can require specific reverse transcriptase-DNA polymerase combinations (see, for example, Chandler et al., 1998). Finally, heat inactivation of the reverse transcriptase before the addition of the DNA polymerase negates the advantages of the coupled RT-PCR and carries with it all the disadvantages of uncoupled RT-PCR systems discussed earlier.\nBecause of the importance of RT-PCR applications, a one-step RT-PCR system with reduced RT inhibition, in the form of a generalized ready-to-use composition, which exhibits high sensitivity, requires a small amount of initial sample, reduces the amount of practitioner manipulation, minimizes the risks of contamination, minimizes the expense of reagents, is not restricted to the use of specific reaction buffers, and maximizes the amount of nucleic acid end product is needed in the art."}
{"text": "1. Field of the Invention\nThe invention relates to a locking detection circuit used in a PLL (phase locked loop) frequency synthesizer for detecting whether the PLL is locked.\n2. Description of the Related Art\nReferring to FIG. 1, a conventional PLL frequency synthesizer 10 comprises a reference frequency demultiplier counter 11, a comparison frequency demultiplier counter 12, a phase comparator 13, a charge pump 14, a low pass filter (hereafter abbreviated as LPF) 15, a voltage controlled oscillator (hereafter abbreviated as VCO) 16 and a lock detection circuit 17.\nThe reference frequency demultiplier counter 11 produces a reference signal fr from a signal of generated by a crystal oscillator 18 through frequency demultiplication. The comparison frequency demultiplier counter 12 produces a compared signal fp obtained from an output signal fv from the VCO 16 through frequency demultiplication. The phase comparator 13 produces a first and a second phase difference signals xcfx86R, xcfx86P in accordance with a phase difference between the reference signal fr and the compared signal fp. On the basis of both phase difference signals xcfx86R, xcfx86P through the operations of the charge pump 14 and the LPF 15, the magnitude of voltage of a control signal VT which is input to the VCO 16 is changed. The PLL circuit 10 also operates to lock the frequency of the output signal fv from the VCO 16 to a desired frequency.\nThe lock detection circuit 17 receives the first and second phase difference signals xcfx86R, xcfx86P from the phase comparator 13, and also receives a reference clock signal CK from the reference frequency demultiplier counter 11 which is obtained by the frequency demultiplication of the signal of from the crystal oscillator 18 at a given ratio. The lock detection circuit 17, which operates in synchronism with the reference clock signal CK, detects whether the output signal fv is locked on the basis of the first and second phase difference signals xcfx86R, xcfx86P, and generates a locking detection signal LD having a level which depends on the result of such detection.\nReferring to FIG. 2, there is shown a specific circuit arrangement of the lock detection circuit 17. As shown, the lock detection circuit 17 includes a NAND circuit 21 that receives the first and the second phase difference signals xcfx86R, xcfx86P from the phase comparator 13 and provides an output signal S1 corresponding to the phase difference which is represented by each pulse width of the signals xcfx86R, xcfx86P. A data flip-flop circuit (hereafter referred to as FF circuit) 22 has a data terminal D for receiving the output signal S1 and a clock terminal CK for receiving the reference clock signal CK, and delivers an output signal S2 corresponding to the output signal S1 at its output terminal Q in synchronism with the rising edge of the reference clock signal CK.\nA NAND circuit 23 receives the signals S1, S2, and delivers its to an inverter circuit 24. An inverted signal S3 is supplied to a data terminal D of an FF circuit 25 from the inverter circuit 24.\nThe FF circuit 25 has a clock terminal for receiving the reference clock signal CK, and provides an output signal S4 at its output terminal Q which depends on the inverted signal S3 in synchronism with the rising edge of the reference clock signal CK.\nAn inverter circuit 30 receives the output signal S4 and generates an inverted signal S4a. A synchronous counter is formed by a plurality of FF circuits 27, 28, 29. The first stage FF circuit 27 has a data terminal D, to which the inverted signal S4a is applied. Each of the FF circuits 27 to 29 has a clock terminal, to which an inverted signal S1a, formed by an inverter circuit 26 with the signal S1, is applied. The FF circuit 27 delivers an output signal S5 at its output terminal Q in synchronism with the rising edge of the inverted signal S1a (or the falling edge of the signal S1). The FF circuit 28 has a data terminal D, to which the output signal S5 is applied, and delivers an output signal S6 at its output terminal Q in synchronism with the falling edge of the output signal S1. The FF circuit 29 has a data terminal D, to which the output signal S6 is applied, and delivers an output signal S7 at its output terminal Q in synchronism with the falling edge of the output signal S1. The output signals S5, S6, and S7 are input to a NAND circuit 31, which then delivers the locking detection signal LD.\nIn the lock detection circuit 17, when one or both of the phase difference signals xcfx86R, xcfx86P has an L level, the NAND circuit 21 delivers the signal S1 which has an H level. The phase difference signals xcfx86R, xcfx86P each have a pulse width which is related to a phase difference between the reference signal fr and the compared signal fp, as will be further described later. Accordingly, the NAND circuit 21 delivers the signal S1 of the H level for a time interval corresponding to the phase difference between the signals fr, fp. The greater the phase difference between the signals fr, fp, the longer the pulse width of the signal S1 or vice versa.\nThe lock detection circuit 17 detects whether the PLL circuit 10 is locked on the basis of the number of rising edges of the reference clock signal CK which are input during a time interval corresponding to the pulse width of the output signal S1 or a time interval during which the NAND circuit 21 delivers the output signal S1 having the H level, and delivers the locking detection signal LD having a level which depends on the result of such detection. Thus it will be seen that the lock detection circuit 17 requires the reference clock signal CK of a higher frequency than the frequencies of the reference signal fr and the compared signal fp. Hence, the reference frequency demultiplier counter 11 produces the reference clock signal CK by the frequency demultiplication at a ratio which is less than the ratio of frequency demultiplication applied to the reference signal fr. Alternatively, the reference frequency demultiplier counter 11 may deliver the input crystal oscillator signal of directly as the reference clock signal CK.\nThe synchronous counter delivers the locking detection signal LD having an H level only when a phase coincidence is reached between the reference signal fr and the compared fp a number of times which is equal to the number of counter stages or more. This prevents the locking detection signal LD having the H level from being delivered from the lock detection circuit 17 for an accidental phase coincidence between the both signals fr, fp.\nDigital mobile equipment generally requires the output signal fv of a higher frequency than analog mobile equipment, and consequently, the PLL circuit 10 produces the reference signal fr and the compared signal fp of higher frequencies, which then approach the frequency of the reference clock signal CK. This may result in a malfunctioning of the lock detection circuit 17.\nFor example, if the PLL circuit is locked between two consecutive rising edges of the reference clock signal CK, the lock detection circuit 17 may be unable to detect the locked condition, thus undesirably delivering the locking detection signal LD having the L level. Because the locking detection signal LD is used in controlling the charge pump 14, the LPF 15 or other external circuit, there are adverse influences upon the operation of the entire PLL circuit or external circuit, causing instability in the operation of the mobile equipment.\nAn object of the present invention is to provide a lock detection circuit and a PLL frequency synthesizer capable of reliably detecting a locked condition.\nTo achieve the above objective, the present invention provides a lock detection circuit for detecting whether a phase of a compared signal is locked with that of a reference signal based on first and second phase difference signals that represents a phase difference between the reference signal having a reference frequency and the compared signal having a preset frequency, the lock detection circuit comprising: a clock generating unit for receiving the first and the second phase difference signals and generating a detecting clock signal in synchronism with one of the first and the second phase difference signals, based on the first and the second phase difference signals; and a lock detecting unit for receiving the first and the second phase difference signals and the detecting clock signal, and detecting whether the phase of the compared signal is locked with the pulse of the reference signal based on the relationship between the detecting clock signal and the phase difference between the first and the second phase difference signals, and generating a lock detecting signal.\nThe present invention further provides a PLL synthesizer comprising: a voltage control oscillator for generating a frequency signal corresponding to a value of a control voltage signal; a comparison frequency demultiplier for generating a compared signal by frequency-demultiplying the frequency signal from the voltage control oscillator; a reference frequency demultiplier for generating a reference signal by frequency-demultiplying an oscillation signal; a phase comparator for receiving the reference signal and the compared signal to compare the phases thereof, and generating first and second phase difference signals, representing a relationship between the reference signal and the compared signal, based on a result of the phase comparison; a charge pump for converting the first and the second phase difference signals from the phase comparator to voltage signal; a low pass filter for receiving the voltage signal from the charge pump and generating the control voltage signal provided to the voltage control oscillator; and a lock detection circuit for detecting whether a phase of the compared signal is locked with a phase of the reference signal based on the first and second phase difference signals and generating a lock detecting signal, the lock detection circuit comprising: a clock generating unit for receiving the first and the second phase difference signals and generating a detecting clock signal in synchronism with one of the first and the second phase difference signals, based on the first and the second phase difference signals; and a lock detecting unit for receiving the first and the second phase difference signals and the detecting clock signal, and detecting whether the phase of the compared signal is locked with the phase of the reference signal based on the relationship between the detecting clock signal and the phase difference between the first and the second phase difference signals, and generating a lock detecting signal.\nThe present invention provides a lock detection circuit for use with a PLL frequency synthesizer for detecting a locked condition of the synthesizer, the synthesizer including a phase comparator which receives a reference signal and a compared signal and generates first and second phase difference signals therefrom, the lock detection circuit comprising: a phase difference detector which receives the first and second phase difference signals and generates a third phase difference signal which depends on a pulse width of each of the first and second phase difference signals; a clock generator circuit which receives the first and second phase difference signals and produces a detection clock signal synchronized with the third phase difference signal; a plurality of delay circuits connected in parallel with each other and in series with the phase difference detector, each of the delay circuits receiving the third phase difference signal and delaying the third phase difference signal by a different delay time; a plurality of switches connected in series with the plurality of delay circuits; a flip-flop circuit having a data input connected to the delay circuits for receiving the third phase difference signal delayed by a predetermined time period, a clock input connected to the clock generator circuit for receiving the detection clock signal, and a data output for providing a status signal, wherein one of the plurality of switches is selectively turned on to delay the third phase difference signal by a selected predetermined time period; and a locking counter connected to a data output of the flip-flop circuit for receiving the status signal and connected to the clock generator circuit for receiving the detection clock signal, the locking counter counting a number of pulses of the detection clock signal while the status signal is at a predetermined level and generating a locking detection signal therefrom, the locking detection signal indicating locked condition of the synthesizer."}
{"text": "Unmanned air vehicles, i.e. vehicles that do not have a physical pilot on board, can be of great interest not just for military missions but also for civil missions. Furthermore, the emergence of new data and image uptake and interpretation systems makes the number of tasks that can be carried out by this type of vehicles increase more and more in both the civil and the military or police scope.\nA fundamental difference between unmanned air vehicles and manned air vehicles is that in the case of the former, in addition to the aerial system, a ground system and data link means necessary for operating the vehicle are also required. Nevertheless, there may be moments or periods which are shorter or longer in time (in some cases, corresponding to most of the duration of the flight) during which the vehicle must function autonomously. Given that during the flight or “mission” in an unmanned air vehicle it is quite probable that unexpected events or conditions (for example changes in weather conditions, wind, turbulences, mechanical problems, etc.) may occur, autonomous control of the flight becomes a complex task.\nThere are a large number of unmanned vehicle control systems. They are usually based on different modules responsible for different parts of the control. For example, the following general modules schematically shown in FIG. 1 may exist:                sensors 101 acquiring and transmitting data related, for example, to the state of the actuators 104, the state of the aircraft (for example, its position, altitude or orientation) and the meteorological conditions (mainly wind intensity and direction);        actuators 104 providing the position of the mechanical control elements which, in the case of an aircraft, provide the forces necessary for controlling the flight;        an estimation and navigation module 102 responsible for obtaining the state variables needed for controlling the system from the values provided by the measurement variables of the sensors;        a guidance and control module 103 providing the actuators 104 with the control variables needed for stabilizing and taking the system state variables to the desired reference values in each case; and        a mission management module 105 which, based on the available data regarding the information on the state of the aircraft it receives from the estimation and navigation module 102 and when the vehicle is flying under the control of an external operator 106, such as a ground control station, based on the instructions it receives from said external operator, provides the guidance and control module with the desired reference variables so as to fulfill certain objectives; this module normally includes means for storing data indicating a mission route comprising a plurality of mission route segments (defined, for example, by “waypoints” corresponding to the mission route).        \nThere are a large number of publications reflecting different aspects of unmanned air vehicle control.\nFor example, patent document U.S. Pat. No. 6,122,572 describes an unmanned air vehicle control system designed for the execution of a mission and having a programmable decision unit capable of managing and controlling the execution of the mission taking into account all available systems and data in the vehicle.\nPatent document U.S. Pat. No. 6,349,258 relates to a method for generating, from two waypoints, a course which must necessarily pass between these two points.\nProgrammed unmanned air vehicles are known to fly according to a “mission route” (which can be preprogrammed) and with the capacity to calculate alternative routes in the case of incidents. For example, patent document U.S. Pat. No. 6,377,875s describes an unmanned air vehicle control system in which a safe flight route is programmed. The vehicle can be controlled remotely via radio; if communication with the control station is lost, the on-board system recalculates the route without the intervention of the control station.\nHowever, the recalculation of the route on board the vehicle requires that the vehicle has an on-board system with sufficient capacity to recalculate the route. This may involve, for example:                the need to have fairly detailed data on the terrain (a digital terrain model);        a complex computer system with the capacity to completely recalculate the route;        a certain risk of “unpredictability” of the route which is finally chosen by the vehicle (which may involve risks and problems for aviation and/or air control in the area, for high buildings in the area, etc.);        uncertainty regarding where recovery of the vehicle will occur;        uncertainty regarding the needs of the vehicle with respect to fuel (given that in the moment of vehicle take-off, its course in the event that it has to divert from its mission route cannot be foreseen).        \nAn objective of the invention is to provide an alternative system for implementing alternative or auxiliary routes apart from the mission route which may involve improvements in some or in all of the aforementioned aspects."}
{"text": "Central/branch retinal artery/vein occlusion, diabetes, glaucoma and, possibly, age related macular degeneration (AMD) are conditions associated with retinal ischaemia. All these diseases may lead to severe sequelae. Therefore, the management of retinal ischaemia is crucial.\nAfter ischaemia/reperfusion (I/R), large amounts of reactive oxygen species (ROS) such as H2O2 are produced. These ROS attack nearby cells and cause tissue damage. Moreover, excessive release of excitatory transmitters such as the glutamate from ischaemia affected neurons leads to neuronal overstimulation and unwanted depolarisation. Consequently, neurons that possess a high density of glutamate receptors are most at risk. This explains why neurons such as retinal ganglion cells (RGCs) and amacrine cells, as well as their neuronal processes, which are located in the inner retina, are vulnerable to (I/R).\nIschaemia induces angiogenesis. Furthermore, in the retina, angiogenesis is often disorganized and typically results in oedema and haemorrhage; these have adverse effects on visual function. There is an urgent need for therapies that promote endogenous protective responses and prevent harmful angiogenesis. Increased levels of hypoxia-inducible factor-1Alpha (HIF-1 Alpha) have been found to be present after retinal ischaemia. HIF-1 binds to the hypoxia response element in hypoxia-responsive target genes, and triggers the expression of vascular endothelium growth factor (VEGF) and matrix metalloproteinases (MMPs). Liu et al. have shown that oxidative stress in human retinal pigment epitheliums (hRPEs) results in up-regulation of VEGF and MMP-9. Additionally, ischaemia has been proved to result in irreversible RGC loss that is accompanied by MMP-9 up-regulation. All the above evidence suggests that the over-expression of HIF-1 Alpha, VEGF and/or MMP-9 in the retina or in RGCs is directly related to ischaemic/ischaemic-like insult, but the relationship in more detail is unknown.\nS-allyl L-cysteine (SAC), an active organosulfur compound in aged garlic extract, has been reported to possess antioxidative activity. In macrophages and endotheliums, SAC has been shown to exhibit potent antioxidative effects involving the scavenging of superoxide radicals, hydroxyl radicals and hydrogen peroxide."}
{"text": "It is desirable to provide decoupling capacitance in a close proximity to an integrated circuit chip or die. The need for such capacitance increases as the switching speed and current requirements of chips or dies becomes higher. Thus, the need for a high number of passive components for high density integrated circuit chips or dies, the resultant increasing circuit density of printed wiring boards (PWB), and a trend to higher frequencies in the multi-gigaHertz range are among the factors combining to increase pressure on passive components surface-mounted on package substrates or PWBs. By incorporating embedded passive components (e.g., capacitors, resistors, inductors) into the package substrate or PWB, improved performance, better reliability, smaller footprint, and lower cost can be achieved.\nCapacitors are the predominant passive component in most circuit designs. Typical materials for suitable embedded capacitor components, such as polymer and high-dielectric constant (high-k) ceramic powder composites or high-k ceramic powder and glass powder mixtures, are generally limited to a capacitance density on the order of nanoFarad/cm2 and 0.1 microFarad/cm2. Attempts have been made to embed thin film capacitors into organic substrates, such as utilizing ceramic fillers in polyimide or epoxy resins in thin laminate form. However, processing and handling of thin-core laminates has proved to be difficult."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nManufacturers are increasingly producing vehicles having higher levels of driving automation. Features such as adaptive cruise control and lateral positioning have become popular and are precursors to greater adoption of fully autonomous-driving-capable vehicles.\nWhile availability of autonomous-driving-capable vehicles is on the rise, users' familiarity and comfort with autonomous-driving functions will not necessarily keep pace. User comfort with the automation is an important aspect in overall technology adoption and user experience.\nDrivers or other vehicle occupants using the autonomous functionality may not be accustomed with various aspects of being a passenger or an autonomous vehicles, lowering satisfaction with the riding experience.\nDrivers using autonomous functionality may also not be accustomed to taking back control of the vehicle, which may be stressful."}
{"text": "When using a heating system to provide ambient air at a comfortable temperature indoors, it is often desirable to increase the humidity of the heated air. If the air to be heated is initially very cold, such as outdoor air in winter, the heated air will contain very little moisture. Indeed, indoor air may have a humidity, or moisture content, as low as 10 to 15 percent, which is significantly drier than even desert air. Such dry air can cause skin and mucous membranes to dry out, making people uncomfortable and perhaps increasing their susceptibility to infection by air-borne viruses and bacteria. Dry air can also lead to drying of wood structures, such as flooring and furniture, and it can lead to increased static electricity that may damage electronic equipment, such as computers. Thus, indoor air is often humidified, using one of several types of humidification systems that are currently available for residential and commercial use.\nIn a bypass type of humidification system, air passes over or through a pad or sheet of wet material that has a large surface area. As the air passes the wet surface, it picks up water that evaporates from the surface. Other systems use nozzles to spray small water droplets into the air. If the droplet size is small enough, the droplets evaporate quickly in the air, thereby preventing significant fogging. These types of systems generally require significant maintenance to prevent scale buildup on the pads or sheets or in the nozzles.\nSteam humidifiers use heat to boil water in a reservoir or tank, most commonly with a heating coil located inside the tank. This type of humidifier generally uses tap water, which may contain significant quantities of dissolved minerals. As the water in the tank evaporates, the mineral content of the remaining water increases, and eventually minerals precipitate from the water to form scale deposits on the walls of the tank and the heating coil. Over time, the tank volume decreases, and the scale layer on the heating coil acts as a thermal insulator. As a result, the humidifier becomes increasingly inefficient. Although frequently flushing clean water through the tank can slow down the rate of scale buildup, there is no simple way to remove the scale deposits once they form in the tanks and in associated plumbing, and the tanks must be replaced periodically.\nThus, there is a need for a steam humidifier that operates more reliably over long periods of time and that can be easily maintained, rather than routinely replacing major system components, such as holding tanks, heating elements, and solenoid valves."}
{"text": "Various systems require electrical coupling between electrical devices disposed within a sealed enclosure or housing and devices or systems external to the enclosure. Oftentimes, such electrical coupling needs to withstand various environmental factors such that a conductive pathway or pathways from the external surface of the enclosure to within the enclosure remains stable. For example, implantable medical devices (IMDs), e.g., cardiac pacemakers, defibrillators, neurostimulators and drug pumps, which include electronic circuitry and one or more power sources, require an enclosure or housing to contain and seal these elements within a body of a patient. Many of these IMDs include one or more electrical feedthrough assemblies to provide electrical connections between the elements contained within the housing and components of the IMD external to the housing, for example, one or more sensors, electrodes, and lead wires mounted on an exterior surface of the housing, or electrical contacts housed within a connector header, which is mounted on the housing to provide coupling for one or more implantable leads, which typically carry one or more electrodes and/or one or more other types of physiological sensors. A physiological sensor, for example a pressure sensor, incorporated within a body of a lead may also require a hermetically-sealed housing to contain electronic circuitry of the sensor and an electrical feedthrough assembly to provide electrical connection between one or more lead wires, which extend within the implantable lead body, and the contained circuitry."}
{"text": "1. Field of the Invention\nThe present invention relates to a cathode-ray tube and a method for manufacturing a cathode-ray tube suitable for use in a projection-type display, for example.\n2. Description of the Related Art\nFor cathode-ray tubes for use in projectors and the like, what is called deposition-type glass (deposited bulbs), in which the panel and funnel sections of the body of the cathode-ray tube are integrally formed, is used.\nBy reason of glass formation, the inner wall of the funnel section and the side of the panel section generally constitute mirror finished surfaces in cathode-ray tubes of this type.\nWhen an electron beam is irradiated upon the fluorescent surface created on the inner surface of the panel section, part of it is reflected off in the form of reflected electrons.\nSince the inner wall of the funnel section and the sides of the panel section constitute mirror-like surfaces, as described above, the reflected electrons are reflected back off the inner wall and are again incident upon the fluorescent surface. This has a marked deteriorating effect on the contrast of the projected image.\nIn the case where a plurality of cathode-ray tubes are located next to each other, such as cathode-ray tubes in a television receiver using a projector, light leaking from the funnel section and the sides of the panel section of a cathode-ray tube affects each other and causes deterioration of the contrast.\nWhen an excellent quality of displayed image is to be achieved in a monochrome cathode-ray tube like those used in projectors, an especially high degree of intensity and contrast is required.\nIn order to counteract factors which cause the contrast to deteriorate, one measure that has conventionally been adopted involves coating the inner wall of the funnel and the sides of the panel with a black material, such as carbon, thereby reducing and absorbing the reflection of light as much as possible.\nWhen the measure is employed for a cathode-ray tube using deposition-type glass in which the panel and funnel sections are integrally formed, coating the inner surface is performed by inserting a brush or other coating implement through the neck section and rotating the glass.\nAccordingly, it has hitherto been possible to coat only that part of the inner surface of the funnel which has a circular cross-section. It also has been difficult to ensure that the funnel section and the sides of the panel section are coated evenly with the black material.\nMeanwhile, if a thick coating is applied to the inner wall in such a manner as to avoid missing anywhere, the black material tends to peel off, thereby causing the withstand voltage properties of the cathode-ray tube to deteriorate.\nThe inner surface of the funnel section needs a furnished conductive film in order to ensure that high voltage supplied from the anode button is fed to the fluorescent surface and the inner surface of the funnel section of the body of the cathode-ray tube has the same potential.\nSince it is possible to coat only that part of the inner surface of the funnel which has a circular cross-section in a cathode-ray tube using deposition-type glass, it has been necessary instead to deposit aluminum over a wide area as a metal backing layer in order to form an inner conductive film.\nHowever, depositing a metal backing layer over a wide area in this manner, in order to form the conductive film, results in the fact that the contrast of the image deteriorates by reason of reflected electrons and stray light caused by not being absorbed sufficiently and reflected off the metal-backing layer.\nExamples of stray light which causes the contrast of the image to deteriorate include stray light caused by light reflected off the inner surface of the body of the cathode-ray tube being subject to multiple reflection within the glass thereof and returning to the front surface of the panel and stray light caused by light reflected off the inner surface of the funnel section of the body of the cathode-ray tube returning directly to the front surface of the panel.\nIn order to solve the abovementioned disadvantages, the present invention proposes a cathode-ray tube having an excellent quality of image with a highly enhanced contrast and a method of manufacturing such a high-contrast cathode-ray tube by reducing stray light and the re-entering of reflected electrons through the fluorescent surface thereof."}
{"text": "Some service providers use conventional risk-based authentication systems to assess risks of processing customer transactions. For example, an online bank may employ a risk engine of such a risk-based authentication system to assign risk scores to banking transactions where higher risk scores indicate higher risk.\nIn generating a risk score, the risk engine takes, as input values, various transaction attributes (e.g., time of receipt, IP address). For each customer of the online bank, there is an associated history based on values of the transaction attributes associated with previous transactions involving that customer. The risk engine incorporates the history associated with the customer into an evaluation of the risk score. Significant variation of one or more attribute values from those in the customer's history may signify that the banking transaction has a high risk.\nFor example, suppose that a particular customer historically submitted transaction requests to the online bank at 3:00 PM from a particular internet service provider (ISP), and, under the customer's identifier, a user submits a new transaction request at 2:00 AM from a different ISP. The different ISP would give rise to a different IP address than that historically associated with the particular customer. In this case, owing to the different IP address and the unusual time that the transaction was submitted, the risk engine would assign a larger risk score to a transaction resulting from the new transaction request."}
{"text": "Lewis Acids have been widely used as catalysts in carbocationic polymerization processes to catalyze the polymerization of monoolefins. Examples of Lewis Acid catalysts include AlCl.sub.3, BF.sub.3, BCl.sub.3, TiCl.sub.4, Al(C.sub.2 H.sub.5).sub.3, Al(C.sub.2 H.sub.5).sub.2 Cl, and Al(C.sub.2 H.sub.5)Cl.sub.2. Such carbocationic polymerization catalysts have many advantages, including high yield, fast reaction rates, good molecular weight control, and utility with a wide variety of monomers. However, conventional carbocationic polymerization processes typically employ Lewis Acid catalysts in unsupported form. Hence, these catalysts, typically, cannot be recycled or reused in a cost effective manner.\nIn a typical carbocationic polymerization process, such as the carbocationic polymerization of isobutylene, a catalyst feedstream in a liquid or gaseous form and a monomer feedstream are fed simultaneously into a conventional reactor. In the reactor, the streams are intermingled and contacted under process conditions such that a desired fraction of the monomer feedstream is polymerized. Then, after an appropriate residence time in the reactor, a discharge stream is withdrawn from the reactor. The discharge stream contains polymer, unreacted monomer and catalyst. In order to recover the polymer, the catalyst and unreacted monomer must be separated from this stream. Typically, there is at least some residue of catalyst in the polymer which cannot be separated. After separation, the catalyst is typically quenched and neutralized. The quenching and neutralization steps tend to generate large quantities of waste which must typically be disposed of as hazardous waste.\nThe recycling or reuse of Lewis Acid catalysts used in polymer processes is difficult because of the chemical and physical characteristics of these catalysts. For example, most Lewis Acid catalysts are non-volatile and cannot be distilled off. Other catalysts are in a solid particulate form and must be separated from the polymer stream by physical separation means. Some Lewis Acid catalysts are gaseous, such as BF.sub.3. The gases can be recycled and reused, but with considerable difficulty, by utilizing gas-liquid separators and compressors.\nThere have been several attempts made to support Lewis Acid catalysts on the surface of inorganic substrates such as silica gel, alumina, and clay. Although these approaches are somewhat successful in recycling the Lewis Acid catalysts, there are several disadvantages associated with their use. One particularly strong disadvantage is that these approaches to supported catalysts generally produce only low molecular weight oligomers. Another disadvantage is that the catalysts (supported on inorganic substrates) typically leach out during the reaction since the catalysts tend to not be firmly fixed to the supporting substrates.\nAttempts to support Lewis Acid catalysts can be characterized as falling into two basic classes; namely, those which rely on physical adsorption and those wherein the Lewis Acid chemically reacts with the support.\nU.S. Pat. No. 3,925,495 discloses a catalyst consisting of graphite having a Lewis Acid intercalated in the lattice thereof.\nU.S. Pat. No. 4,112,011 discloses a catalyst comprising gallium compounds on a suitable support such as aluminas, silicas and silica aluminas.\nU.S. Pat. No. 4,235,756 discloses a catalyst comprising porous gamma alumina impregnated with an aluminum hydride.\nU.S. Pat. No. 4,288,449 discloses chloride alumina catalysts.\nU.S. Pat. Nos. 4,734,472 and 4,751,276 disclose a method for preparing functionalized (e.g., hydroxy functionalized) alpha-olefin polymers and copolymers derived from a borane containing intermediate.\nU.S. Pat. No. 4,167,616 discloses polymerization with diborane adducts or oligomers of boron-containing monomers.\nU.S. Pat. No. 4,698,403 discloses a process for the preparation of ethylene copolymers in the presence of selected nickel-containing catalysts.\nU.S. Pat. No. 4,638,092 discloses organo-boron compounds with strong aerobic initiator action to start polymerizations.\nU.S. Pat. No. 4,342,849 discloses novel telechelic polymers formed by hydroborating diolefins to polyboranes and oxidizing the polymeric boranes to form the telechelic dehydroxy polymer. No use of the resulting polymer to support Lewis Acid catalysts is disclosed.\nU.S. Pat. No. 4,558,170 discloses a continuous cationic polymerization process wherein a cocatalyst is mixed with a monomer feedstream prior to introduction of the feedstream to a reactor containing a Lewis Acid catalyst.\nU.S. Pat. Nos. 4,719,190, 4,798,190 and 4,929,800 disclose hydrocarbon conversion and polymerization catalysts prepared by reacting a solid adsorbent containing surface hydroxyl groups with certain Lewis Acid catalysts in halogenated solvent. The only disclosed adsorbents are inorganic; namely, silica alumina, boron oxide, zeolite, magnesia and titania.\nU.S. Pat. No. 4,605,808 discloses a process for producing polyisobutene using a complex of boron trifluoride and alcohol as catalyst.\nU.S. Pat. No. 4,139,417, discloses amorphous copolymers of monoolefins or of monoolefins and nonconjugated dienes with unsaturated derivatives of imides. In the preparation of the polymer the imide is complexed with a Lewis Acid catalyst.\nJapanese Patent Application No. 188996/1952 (Laid Open No. J59080413A/1984) discloses a process for preparing a copolymer of an olefin and a polar vinyl monomer which comprises copolymerizing an olefin with a complex of the polar vinyl monomer and a Lewis acid.\nEuropean Patent Application No. 87311534.9 (Publication No. EPA 0274912) discloses polyalcohol copolymers made using borane chemistry.\nT. C. Chung and D. Rhubright, Macromolecules, Vol. 24, 970-972, (1991) discloses functionalized polypropylene copolymers made using borane chemistry.\nT. C. Chung, Journal of Inorganic and Organometallic Polymers, Vol. 1, No. 1, 37-51, (1991) discloses the preparation of polyboranes and borane monomers.\nU.S. Pat. No. 4,849,572 discloses a process for preparing polybutenes having enhanced reactivity using a BF.sub.3 catalyst. Polybutene is produced which has a number average molecular weight in the range of from 500 to 5,000. The polymer has a total terminal double-bond content of at least 40% based on total theoretical unsaturation of the polybutene. The polybutene contains at least 50% by weight isobutylene units based on the polybutene number average molecular weight. The process is accomplished by contacting a feed supply comprising at least 10% by weight isobutylene based on the weight of the feed with a BF.sub.3 catalyst under conditions to cationically polymerize the feed in liquid phase to form polybutene. The polymer is immediately quenched with a quench medium sufficient to deactivate the BF.sub.3 catalyst.\nThere has been a continuous search for catalysts having high efficiency which can be recycled or reused in cationic polymerization processes. The present invention was developed pursuant to this search."}
{"text": "Casual photographers often compose scenes in a manner that is appealing to them when seen through a camera viewfinder, but is later found to be unappealing when seen in a resulting photographic print or other final image. Instruction on how to take better photographs is readily available in books, in classes, and the like; but such modes of instruction are burdensome to access during picture taking and are not much used by casual photographers. This is unfortunate, since the result is that many people repeat the same mistakes over and over, and also miss out on the fun of learning how to take better photographs.\nCameras are known that provide an indicator or lock up the shutter release when a forthcoming shot would be too dark or too close. Verifying cameras provide a verification image to users immediately after capture of a scene on photographic film or other archival media. The verification image is provided on a digital display and portrays the image captured on the archival media. This allows the user to review the verification image and decide if the shot was unsuccessful and should be repeated. Major capture failures, such as a something blocking the lens system, are readily apparent. Other characteristics of the captured image, such as composition, are also shown; but, in view of the small size of the digital display may not be immediately apparent to the user.\nJapanese published patent application No. 07-319034 discloses a hybrid camera in which the photographer can change exposure settings to modify a verification image. The photographer then knows whether to repeat the shot with the changed settings.\nU.S. Pat. No. 5,640,628 discloses a camera that can change metadata indicating a default number of prints, in response to a determined condition.\nHybrid cameras are known that use an electronic image capture unit having a larger angle of view than a corresponding film image capture unit. U.S. Pat. No. 4,978,983 discloses a camera that uses the larger area of the electronic capture unit to correct for parallax at some focusing distances. A display on the camera shows a digital image that corresponds to the angle of view of the film image capture unit.\nSoftware is widely available that allows for the easy manipulation of digital images. Digital cameras can be used to capture images which are then modified using such software after downloading to a computer. This is a powerful approach, but lacks immediacy, since the images are not manipulated on the camera.\nDigital cameras necessarily make some modifications of captured images. Some cameras also allow the user to selectively modify some images. For example, the use of digital zoom is disclosed in U.S. Pat. No. 5,172,234. A problem with these approaches is complexity or lack of immediacy or both. A novice is likely to be confused if he or she attempts to modify images during a picture taking session. Modifying images on a camera after a picture taking session is less confusing, but remains complex unless user choices are strictly limited.\nOne reason for modifying captured images is correction of mistakes by the user and improvement of photographic technique. The widely available educational materials for this purpose are of little help to a user during a picture taking session.\nIt would thus be desirable to provide an improved camera and method in which the camera displays one or more different suggestions on how to improve on a just captured archival image in a later recapture and the camera sets up in configuration for that improved recapture when one of the suggestions is selected."}
{"text": "The present invention relates to a method for tracking a maneuvering target with a slow scan rate sensor, and, more particularly, to tracking a highly maneuvering target without overburdening the target processing capabilities of a tracking system using such a sensor.\nBy \"slow scan rate\" is meant a scan rate, wherein a target monitoring or sensor system, such as a radar antenna or sonar hydrophone, examines a predetermined volume in space and returns to the same volume relative to the sensor in a time interval which is long relative to the anticipated or maximum velocity acceleration and maneuvering capability of a potential target. By way of example and not of limitation, a known radar system employs an antenna which is rotated or slewed through 360.degree. in about ten seconds. Thus for a stationary target, the sensor would return to the subject detection volume nominally every 10 seconds depending on the sensor platform or own ship motion, if any.\nFor certain current jet aircraft or missiles which may be desired to be detected and tracked, an interval as long as ten seconds between obtaining updated actual target position information may permit the target to maneuver far enough away from its previous historical position and path so that a next interval detection using a conventional system cannot be reliably correlated or determined to be the same target.\nTracking targets through their maneuvers is a classical tracking problem solved by conventional tracking algorithms. Successful target track is typically the result of frequent updates of target position during any maneuver. Slow scanning sensors do not provide the required frequent target updates for conventional trackers to maintain track on maneuvering targets. Within one scan interval of a slow scanning sensor, a maneuvering target, like a cruise missile, can initiate, perform, and terminate a large acceleration or high g turning maneuver away from its previous position and course. The term g is used to refer to acceleration in terms of a multiple of the acceleration due to gravity which is nominally 9.8 m/sec.sup.2 On the next scan, the target location may be a significant distance from the location predicted by the tracker from historical data of a predetermined number of prior scans and the target may be headed on a drastically different course from the one derived from such historical data. Either or both of these events may defeat the performance of conventional target correlation and tracking algorithms when employed in conjunction with a slow scanning sensor.\nPreviously developed scan-to-scan correlation methods for use with slow scanning sensors have exhibited a number of deficiencies such as excessive false alarms, target tracking limitations and/or intermittent target tracking. Some slow scan tracking methods employ small correlation gates to restrict the probability of false correlation,thereby limiting the false track report rate. These trackers do not have the capability to track maneuvering targets because of the insufficient size of the gates. Other tracker developments for use with such sensors have utilized large, non-optimally shaped correlation gates which can track maneuvering targets but which permit an excessive number of false correlations and false target reports. Further, use of such large gates inefficiently consumes target processing resources of the overall system, such as a radar or sonar by examining portions of a spatial volume where there is a low probability of finding the historically tracked target."}
{"text": "The present invention relates to a portable terminal, an information processing apparatus, a content display system, and a content display method.\nCellular phones capable of browsing internet sites and television (TV) programs have come into wide use today. On the other hand, there have been increasingly employed television sets capable of connecting to networks. In this environment, a technique in which a content received by a cellular phone is viewed by a television set with higher usability and a technique in which a cellular phone is employed as a remote control device to operate a television set are under discussion.\nFor example, JP-A-2006-286855 describes a cellular phone including a broadcast receiving function. The cellular phone transmits a history of viewed programs and playback indication signals thereof to a video playback device. According to the cellular phone, a program which the user viewed outside his or her house can be again viewed on a large screen after returning home. It is also possible that when the user viewed part a program outside home during a free period of time, the user can view its subsequent part, parts before and after the viewed part, or the entire program from the start point to the end point thereof after returning home by using a recording and reproducing device in the house,\nJP-A-2005-135346 describes a technique for use in a system including a cellular phone, an information processing apparatus, and a display. When mail data including Uniform Resource Locator (URL) information is received, the cellular phone transfers the mail data via radio communication to the information processing apparatus. On receiving the mail data, the information processing apparatus displays the mail data on the display. If an operation is conducted for the URL information in the mail data, an operation is conducted to access a radio communication network other than a network used for the radio communication on the basis of the URL information, and resultantly obtained information is displayed on the display.\nJP-A-2006-338406 describes a technique in which a communication terminal device including another function is ordinary employed as a remote control device to remotely control other communication terminal devices.\nAdditionally, JP-A-2007-74265 describes a videoconference system using a cellular phone connected via the internet to a videoconference server to communicate information with the server and a display device capable of, when connected via a communication function to the cellular phone, displaying the information. The videoconference system controls devices connected thereto via the videoconference server and a network. The videoconference system includes an information display module to display on a first display of the cellular phone a first videophone screen on which the information is displayed using letters. The information display module also displays on a second display of the display device a second videophone screen on which image information included in the information communicated as above is displayed. JP-A-2007-74265 describes a system in which a monitor camera is remotely controlled by a cellular phone, and a video image produced from the camera is displayed on a display screen of a Personal Computer (PC)."}
{"text": "Some solutions have a relatively short useful life after a solute is mixed with a solvent. Once such solution is a thrombolytic agent which is supplied intravenously to a patient to dissolve blood clots. The solution must be used within a few hours after the thrombolytic agent, in powder form, is dissolved in a liquid carrier such as water. At the present time, when the solution is to be administered, a nurse or other person adds a prescribed quantity of the liquid solvent into a vessel containing a prescribed quantity of the thrombolytic powder and shakes or otherwise mixes the mixture to accelerate the dissolution of the powder in the solvent. The solution must be administered shortly thereafter.\nIt would be desirable to provide a more sterile and foolproof system for shipping, storing and subsequently mixing a solvent and a solute to insure a precise mixture of premeasured quantities of the solvent and the solute under completely sterile conditions."}
{"text": "1. Field of the Invention\nThe present invention relates to a hydraulic circuit, and in particular to a hydraulic circuit including an accumulator having an inflow passage which introduces a hydraulic fluid which is discharged from a hydraulic pump into a hydraulic fluid chamber and a discharge passage which discharges the hydraulic fluid from the hydraulic fluid chamber to a hydraulic actuator.\n2. Description of the Related Art\nAn example of this type of hydraulic circuit is disclosed in, for example, Japanese Patent No. 2576998. In the disclosed hydraulic circuit, a hydraulic fluid which is discharged from a hydraulic pump is introduced into a hydraulic fluid chamber of an accumulator through an inflow passage, and then the hydraulic fluid is discharged from the hydraulic fluid chamber of the accumulator to a hydraulic actuator such as a hydraulic booster through a discharge passage. As a result, pulsations of the hydraulic fluid discharged from the hydraulic pump are securely decreased by the operation of the accumulator.\nIn the case where an accumulator which operates when the pressure in a hydraulic fluid chamber is at least a set pressure is used to decrease pulsations of the hydraulic fluid discharged from a hydraulic pump, in a transient period until the pressure in the hydraulic fluid chamber reaches the set pressure, the accumulator does not operate, and pulsations of the hydraulic fluid which is discharged from the hydraulic pump cannot be decreased."}
{"text": "Various applications of fluorescence techniques to analyze biological samples are known to people skilled in the art. In case of electrophoretic techniques proteins or DNA are labeled with a fluorescence probe to visualize their electrophoretic bands in gels or columns. In addition, most biochip applications so far are based on a fluorescence read-out, whereas the specific binding of a fluorescence-labeled target molecule to a probe molecule immobilized on a solid support is monitored. Applications for DNA analysis in the liquid phase include fluorescence hybridization probes like the double-stranded DNA binding dye SybrGreenI or FRET (Fluorescent Resonance Energy Transfer) probes utilizing two fluorescence probes and energy transfer. A very important application for fluorescence techniques in the liquid phase is the quantification of PCR products in real time, the so-called real-time PCR.\nIn all these cases, a fluorescence reading device is needed that provides light of a certain wave length to excite the fluorescence label of the assay and that is able to detect the fluorescence light form said label emitted at a somewhat different wavelength. One major problem of all fluorescence reading devices is the enormous intensity of the excitation light in comparison with the fluorescence light emitted by the dye and therefore, one has to assure that the excitation beam does not hit the detector in order to monitor the fluorescence signals accurately. In other words, the optical path of the excitation light has to be different from the optical path of the fluorescence light, at least partially.\nThe realization of the fluorescence principle is quiet easy, when only one fluorescence probe has to be monitored in the liquid phase of e.g. a capillary. Here, e.g. a white light source together with a set of dichroic mirrors and filters is sufficient to meet the requirements. However, if more than one fluorescence label is present in the sample, a lateral distribution of spots on a solid support or the fluorescence of a microtiter plate has to be monitored, the requirements for the fluorescence reading device are more difficult to fulfill.\nIn principle, there are two different strategies to excite and monitor the fluorescence of a lateral distribution of sites. The first strategy is to scan the lateral distribution of sites, whereby the individual sites are successively analyzed one at a time. The second strategy is to illuminate the whole distribution of sites simultaneously and to image the corresponding fluorescence e.g. on a CCD chip. The scanning strategy has the obvious drawback that either the support has to be moved in two dimensions (WO 03/069391, DE 102 00 499), the detector has to be moved with respect to the support (US 2002/159057), the detector has to move in one dimension and the support in the other dimension or the optics has to include one or two dimensional scanning means i.e. galvo mirrors. On the other hand, the main difficulty of the strategy to illuminate the whole support simultaneously is to assure a homogeneous illumination across the whole distribution of sites. An alternative to the homogeneous illumination of the whole distribution of sites is the use of an array of light sources, whereby each site is illuminated by its own light source. DE 101 31 687 describes this strategy for the evaluation of PCR in a thermocycler with a plurality of wells using a beam splitter and an array of LEDs for illumination. DE 101 55 142 describes the dark field monitoring of fluorescence signals, wherein the microarray is illuminated by an array of LEDs, too, but in this embodiment no beam splitter is needed.\nConcerning the requirement to separate the optical path of the excitation beam and of the fluorescence light at least partially, there are again two different possibilities. The first possibility is the so called epi-illumination, whereby beam splitters are utilized and the excitation beam and the fluorescence light share at least part of the optical train. The second possibility is the use of oblique illumination. Here, the excitation beam is arranged in such a way that it has a certain angle to the normal of the support surface and the corresponding reflection of the excitation beam is outside of the acceptance angle of the detection system (e.g. US 2002/0005493 A1, EP 1 275 954 A2).\nUS 2003/0011772 A1 describes an optical apparatus to simultaneously observe a plurality of fluorescence dyes in a probe using a beam splitter. DE 197 48 211 A1 discloses a system to monitor the fluorescence signals generated in the wells of a microtiter plate simultaneously using a beam splitter, a field lens and an array of lenses focusing the light into each well. The detection is performed by imaging the light onto an array of photodiodes or a CCD chip. The fluorescence light collected in this embodiment of the system is appointed by the amount of dyes excited by the light cone of the focusing lens and therefore is dependent on the fill level of the well. WO 99/60381 claims an instrument for monitoring PCR reactions simultaneously in a plurality of vials in a temperature cycled block. The optical components of this instrument include again a beam splitter, a field lens, an array of vial lenses focusing individual light beams into each vial and a detection mean focusing the emission light onto e.g. a CCD detector. Due to the necessity of an array of vial lenses, the size and the lateral density of individual sites is limited. The JP 2002014044 describes a fluorometric apparatus to monitor fluorescence generated at a plurality of wells. The optical components comprise a beam splitter and a lens system to illuminate the wells collectively with light being parallel to the direction of the depth of the wells. However, the image forming optical system condenses the light onto a detection mean. U.S. Pat. No. 6,498,690 B1 discloses a method for imaging assays with an objective comprising a telecentric lens. U.S. Pat. No. 6,246,525 B1 claims an imaging device for imaging a sample carrier comprising a Fresnel lens.\nThus, it was the object of the present invention to provide an improved device for simultaneous monitoring of fluorescence signals from a lateral distribution of sites by optimizing the optical path towards homogeneous illumination and accurate detection. In one aspect of the present invention, the problem to be solved relates to improvements in monitoring multiplexed real-time PCR in a microtiter plate format."}
{"text": "A thermoelectric cooler (TEC), also known as a Peltier cooler, is a solid-state electrical device that may be configured to transport heat when current is passed through a number of semiconducting “pellets.” The pellets are typically configured in a series circuit arranged to produce a desired degree of cooling and device resistance. The direction of heat transport in a TEC may be determined by the direction of current flow through the pellets.\nTECs provide a convenient and effective means of temperature control in many applications. In one such application, these devices have been used in electronics systems to reduce the operating temperature of electronic components. Such cooling is especially appropriate where system design constraints preclude the use of cooling fins or forced air flow. TECs may also be used to refrigerate a component by cooling the component below the ambient temperature. Also, by coupling the TEC to a feedback system, a TEC may be used to regulate the temperature of a device by operating in both a heating and a cooling mode."}
{"text": "This invention relates to a mattress which provides a sound sleep to a sleeper.\nConventionally, it has been believed that bedding or mattress having a soft feeling assures a sound sleep. Furthermore, with the advent of foam rubber and spring mattresses, the type of bedding used today feels much softer.\nHowever, these types, in general, are so soft that the spinal cord of a sleeper tends to bend downwardly during sleep. Therefore the weight of the sleeper cannot be uniformly supported by the bedding.\nDue to the above manner of sleeping, a person often feels pain in his shoulder, waist or spinal cord.\nTo be more specific, since conventional bedding is made of extremely soft material, such as urethane foam rubber, the central portion (e.g. shoulder, waist or hip) of the sleeper's body sinks into the bedding.\nTherefore, sufficient respiration to impart a sound sleep cannot be obtained while the blood is congested and the body temperature is high in those areas which come into direct contact with the mattress. Both of these prevent the sound sleep which also occurs on the part of the body (e.g. shoulder, hip or waist) that sinks into the bedding.\nOn the other hand, some people believe that a bedding made of hard material assures a sound sleep. However, in practice, this type of bedding is not suitable for any person except youngsters, since it causes the upward curving of the spine and the stretching of muscles.\nTherefore, it is of vital importance that the ideal condition of a mattress be such that it is soft to the touch yet firm enough to uniformly distribute the sleeper's weight.\nIt is an object of the present invention to provide a mattress which can resolve the afore-mentioned defects of conventional beddings, including the spring mattresses.\nIt is another object of the present invention to provide a mattress which can enhance the circulation of the blood within the body as well as set the sleeper free from pains that occur in the shoulders or areas which come in direct contact with the upper surface of the bleeding.\nIt is still another object of the present invention to provide a mattress which is further equipped with a desired number of magnets which can effectively cure the deficiency syndrome of magnetic fields, such as the imbalance of autonomic nerve or the stiffness in the shoulders."}
{"text": "1. Field of the Invention\nThis invention relates to the general field of optical instrumentation and, in particular, to methods and mechanisms for adjusting the position of the focal point of optics in the optical path of imaging systems.\n2. Description of the Prior Art\nOptical telescopes and their basic components (such as the objective, the eye-piece, and the telescopic optical-tube assembly) are some of the oldest types of optical instrumentation known in the art. They are fabricated by very diverse manufacturers and are often built according to custom designs. As a result of this diversity of fabrication, some opto-mechanical standards in the telescope industry have not been established or strictly maintained.\nAs an illustration, a typical eye-piece 10, shown schematically in FIG. 1, consists of an eye-piece lens 12 mounted inside an eye-piece tube 14. Due to the fact that different manufacturers use eye-piece tubes of various lengths and mount the lenses in different positions inside these tubes, the distance 1 between the focal point FE of the lens 12 and the edge 16 of the tube in otherwise optically equivalent eye-pieces differs from manufacturer to manufacturer. Similarly, as shown in FIG. 2, there is no established standard for positioning the lens 18 of an objective 20 within its housing 21 (known as the objective cell), nor are there standard objective cells. As a result, the exact position of the focal point FO of the objective lens 18 with respect to its housing 21 also differs from case to case.\nTherefore, when a typical telescope 30 is assembled conventionally by combining an eye-piece 10 with an objective 20, the focal points FO, FE of the objective lens 18 and of an eye-piece lens 12, respectively, are invariably separated from each other by a distance d, as illustrated in FIG. 2. For proper operation of the telescope 30, it is required that the foci FO, FE of the coaxial objective anti eye-piece lenses 18, 12 coincide in space. Thus, to compensate for the distance d separating the foci (i.e., to re-focus the telescope), a conventional focusing mechanism is used to mechanically translate the eye-piece and/or the objective lenses with respect to each other along the axis 32 of the telescope. Such a focusing mechanism 33 is typically attached externally to the optical tube assembly 34 of the telescope.\nFor the purposes of this invention, the optical tube assembly (OTA) of a telescope is the portion of the telescope housing connecting the objective cell to the eye-piece. Often additional optical elements are located inside the OTA along the optical train prior to the eye-piece. Conventional focusing mechanisms (not shown in FIG. 2) typically operate as mechanically driven telescope tubes (either threaded, sliding, or geared) adapted to change the dimensions of the telescope (its so-called “foot-print”) by either extending or contracting the OTA.\nThe optical path of the telescope between the objective lens and its focus is often appropriately folded to reduce the overall dimensions of the instrument. This is typically achieved by positioning secondary mirrors between the objective and the eye-piece at an angle designed to reflect the beam away from its initial direction of propagation. As a result, the length of the OTA is reduced at the expense of increasing its lateral dimension. The need for a focusing mechanism, however, exists regardless of whether the imaging system of the telescope is folded or linear. Therefore, when using a conventional focusing mechanism, it is not possible to have a working telescope of fixed dimensions (foot-print), which represents a problem for miniaturization purposes or when the telescope is part of a larger opto-mechanical system that is subject to dimensional constraints.\nTo the extent that mere mechanical adjustment of the length of any part of the telescopic housing is used for re-focusing the telescope, this limitation is unavoidable. Thus, there remains a need for a focusing mechanism that allows the manufacture of a telescope of fixed dimension, especially in miniaturized implementations."}
{"text": "This invention relates to an improved picture hanger. More particularly, it relates to an assembly which may be either integral with or attached to a frame during its manufacture to provide both a laterally adjustable, rugged hanger and a decorative backing for the frame. An example of a picture hanger which is somewhat related to the invention herein disclosed may be found in U.S. Pat. No. 2,875,542 to J. R. Peach."}
{"text": "Despite current social and economic advances, there still remains an overabundance of criminal activity in most societies. Unfortunately, the amount of resources for combating such criminal activity is typically limited. Accordingly, when criminal activity occurs, police and other law enforcement agencies typically rely on citizen calling for aid via specific, well-established communications channels. For example, in the United States, the telephone number 911 is utilized almost universally as a means to contact law enforcement agencies and other emergency response services. However, 911-based systems are limited in some aspects.\nFirst, even though 911 services are accessible via virtually any telephone, they still rely on a dialing process. On a landline telephone or a basic mobile phone (i.e., other than a smartphone) the dialing process is trivial. However, on various other types of devices (smartphones, computers, tablet devices, etc. . . . ), additional steps can be required to establishing the call. For example, on a smartphone, a screen unlocking process may be required, followed by selection of an phone application. Still another difficulty is that in some scenarios, even the simple act of dialing 911 can be difficult. For example, an individual being robbed or assaulted may be under high stress conditions and therefore may be unable to properly dial 911 or operate his phone at all in a normal fashion. In another example, an injured individual may be physically unable to dial 911.\nSecond, seeking assistance via 911 can sometime be an inefficient process. For example, in some circumstances, the closest or most appropriate emergency response service may not be reached by the call. For example, a college or university may have their own on-campus emergency response services. However, a mobile phone call to 911 from a person on the campus would not typically reach such services. Rather, the call would go to a central 911 hub from which the campus emergency response services may be dispatched. As a result, it is possible that life threatening delays can occur while the campus emergency response services are contacted or dispatched.\nOne solution to the above-mentioned problems is the use of specialized emergency devices. However, such devices are typically limited. For example, some devices are merely act as a distress beacon, which may communicate location information, but no information regarding the present state of the user. For example, it is not uncommon to see the use of call stations on university campuses that provide one-button access to on-campus emergency response services. Thus, location information is provided when such a station is activated, but no information about the victim is provided. Further, such stations are generally fixed, so they rely on the victim being able to reach the station and for the victim to remain in proximity of the station.\nOther devices may be associated with specific users, but such devices are restricted to sending a signal that contains a limited amount of vital content about the user or may rely on a third party dispatcher. For example, emergency response services such as the LIFE ALERT system of Life Alert Emergency Response, Inc., provide a remote control device that allows users to cause a speakerphone in their home to automatically place a call to a monitoring service. The monitoring service then contacts appropriate persons for responding to the user's emergency. Unfortunately, such systems typically require the use of specialized equipment. Further, such systems typically lack the capability to directly contact emergency services to begin dispatch of emergency response services immediately. As noted above, delays in dispatching emergency response personnel can have adverse, if not fatal, consequences for the user."}
{"text": "Organic electroluminescent devices (OLEDs) comprise very thin layers of organic substances on top of a glass substrate covered with an electrically conducting but optically transparent oxide, usually Indium-Tin oxide (ITO). ITO forms the substrate electrode, usually the anode and a layer of Aluminum (100 nm) evaporated on top of the organic layer stack forms the counter electrode, usually the cathode. When a voltage between 2 and typically 5 V is applied between the electrodes, current is injected into the organic layers and light is generated. The preparation of OLEDs requires structuring of all individual layers to provide a reliable operation over time, in particular to operate the OLEDs, both electrodes have to be electrically isolated from each other. Therefore at least three masks processes are commonly required: structuring of the substrate electrode deposited on top of the substrate (first mask process), depositing the organic layer stack on top of the pre-structured substrate electrode (second mask process different to the first mask process) and depositing the counter electrode on top of the organic layer stack without providing an electrical contact to the substrate electrode (third mask process different to first and second mask process). The mask processes apply so-called shadow masks having different geometries in order to cover the desired areas with the layer to be deposited and simultaneously shielding the other areas in order to avoid material deposition there. The application of shadow masks for deposition (coating) processes is expensive, because the masks have to be manufactured with high geometric accuracy. Furthermore masks will be coated with the deposited material and have to be cleaned periodically. Misalignments of masks may lead to deposition failures making the operation of OLEDs impossible, e.g. by electrically shortened electrodes due to a misalignment of the counter electrode mask. A lot of additional measures have to be applied in order to achieve a reliable coating process making this process very expensive, see for example EP 1202329 A2. All these measures result in a manufacturing process requiring an enormous effort to achieve a good production yield (low failure ratio)."}
{"text": "This invention relates to an evacuated container and method for use in analyzing a centrifuged anticoagulated whole blood sample for the presence or absence of rare events. More particularly, this invention relates to the detection of hematological rare events such as circulating cancer cells, bacteria, hemato-parasites, or other rarely occurring visually or photometrically detectable particles in the blood sample. A preferred embodiment of the container of this invention is an evacuated rectilinear container which contains an insert which container is essentially the same length and width as a microscope slide.\nIt is known that rare events, such as cancer cells may be present in the blood stream. Co-pending U.S. patent application Ser. No. 08/976,886, filed Nov. 24, 1997 describes a method for examining a centrifuged sample of anticoagulated whole blood for the presence or absence of cancer cells or other rare events. The method described in the aforesaid patent application involves the use of a capillary tube containing a cylindrical insert or float which restricts the available space in the tube into which the white cells and platelets will settle during centrifugation. This physically elongates the buffy coat portion of the blood sample and forces any cancer cells in the blood sample toward the tube wall where they can be seen under suitable magnification. We have discovered that any cancer cells present in the blood sample will settle into the area occupied by the buffy coat, and particularly in a region above the granulocytes and below the plasma, during centrifugation. The exact location of the cancer cells in the buffy coat will be governed by the density of the cancer cells and other less well-characterized physical forces.\nThe capillary tubes utilized in connection with the procedure described in the aforesaid patent application are capable of holding about 111 micro liters of blood. The frequency of cancer cells in peripheral blood can be as low as about one cancer cell per milliliter, and will be relatively dependent on the stage of the cancer in the patient being tested. Stated another way, a patient in an early stage of cancer is likely to have a much lower number of circulating cancer cells in an individual\"\"s peripheral blood than a patient in advanced stages of the disease, and a sample of less than about 111 xcexcl of blood stands about a 90% chance of missing a rare event in the blood sample when the rare event occurs at a frequency of about 1 per ml of sample. The ability to examine larger samples of peripheral blood could enable earlier detection of the disease or other rare events. A larger tube and insert could be used to enable the examination of larger blood samples, but tubes produce spherical aberrations due to the curvature of the glass tube wall. It is also noted that the use of tubular containers requires that the container be examined throughout a 360xc2x0 arc for the presence or absence of the rare events. This requires rotating the tube in the microscopical examining instrument.\nIt would be highly desirable to be able to examine larger blood samples in the same manner as provided by the capillary tubes and inserts, without encountering visual aberrations and without the need to rotate the sample chamber, and in which a standard microscope stage could be used without significant modification. However, tubular specimen sample containers can also be used in the performance of the method aspects of this invention.\nThis invention relates to an assembly for analyzing larger samples of anticoagulated whole blood for the presence or absence of circulating cancer cells or other rare events. The assembly comprises a hollow sample container which has essentially the same shape, length and width of a microscope slide, but which is thicker than a microscope slide. The hollow interior of the container can hold from about five to about ten milliliters of blood. One end of the container is sealed and the other end is closed by an elastomeric stopper. Obviously, both ends of the container could be closed by elastomeric stoppers if so desired. A relatively flattened rectilinear insert is disposed in the container and occupies between about 80-95%, and preferably about 90% of the volume of the portion of the interior of the container in the area occupied by the insert. The specific gravity or density of the insert is such that the insert will settle into or float in the packed red cell layer and be surrounded by the buffy coat constituents during centrifugation of the blood sample in the container. The interior of the container is evacuated so that blood will flow into the container as the result of a cannula puncturing the elastomeric closure. The container can be used to draw blood directly from a patient, or can be used to obtain the blood from a larger container such as a VACUTAINER(copyright) brand container sold by Becton Dickinson and Company. When the container assembly of this invention is used to draw blood directly from a patient, the testing reagents needed to anticoagulate the blood sample and detect the rare events in the blood sample can be pre-incorporated into the sample container assembly.\nReagents needed to identify rare events can include stains such as acridine orange which can highlight cell morphology; antibodies which are specific to surface receptors on cancer cells or other rare events; and stains which, when irradiated by light of appropriate wavelengths, will emit signals at wavelengths which can be photometrically and/or visually differentially detected over the background noise of emissions from the blood sample constituents such as white cells, cytoplasm, hemoglobin, and the like, which result when the blood sample is irradiated by such light sources.\nThe insert can include a structure which will skew the location of the insert inside of the container so that the majority of the white cell and platelet constituents will be located on one side of the rectilinear insert thereby enabling that portion of the buffy coat to be examined for the presence or absence of cancer cells, or other rare events without the need of rotating the container about its axis. Suitable skewing structures can include rails, hemispherical bumps, fins, or the like, on the insert or on the interior of the container. The container assembly is preferably sized so as to be positionable in a conventional centrifuge, and also positionable on a conventional microscope stage.\nIt is therefore an object of this invention to provide an evacuated anticoagulated whole blood sampling container assembly, which can be centrifuged, which container assembly is rectilinear in configuration, and which container assembly includes a volume-occupying insert that restricts the volume available in the container assembly wherein rare events will settle during centrifugation of the blood sample in the container assembly.\nIt is a further object of this invention to provide a method for examining a centrifuged anticoagulated whole blood sample under suitable sample illumination conditions that will reveal the presence or absence of rare events, such as cancer cells in predetermined areas of the blood sample."}
{"text": "This invention relates, but is not limited, to optical character recognition systems such as optical character readers and, more particularly, to methods and apparatus for correcting for skew conditions in text lines being optically read.\nA basic requirement for Optical Character Recognition (OCR) is to locate and bound each line of text on a document. For the usual OCR application, this is a constrained task in that apriori it is known that the page contains one font, all character spaces are defined, and all text line spaces are given. Furthermore, the document usually contains only textual material, or, if non-textual data exists, it is field-formatted to eliminate it from consideration in line finding. Text line tilt is carefully controlled, and many commercial OCR machines utilize the top line of the text on a page for a reference to physically deskew the page prior to line finding.\nFor many standard fonts, FIG. 1 is illustrative of the basic types of character elements. Generally, capital letters, such as the \"A\" and \"B\", are bounded by levels 2 and 4. Lower case such as the \"a\", is bounded by levels 3 and 4, and is shorter in stature than the capitals. Lower case such as the \"b\", is constrained by levels 1 and 4, and is taller than the capitals. Lower case such as the \"g\", is bounded by levels 3 and 5, and is about as tall as the capitals, but with a segment dropping below the capital bottoms. Punctuation marks such as the period, quotations, and commas are short in stature, and are attached to either levels 2 or 4. Arithmetic operators such as the plus sign or asterisk are also short, and congregate around level 3.\nA typical situation of skewing is shown in FIG. 2. The document 10 contains lines of text generally indicated as 12. As document 10 moves through an optical character reader in the direction of arrow 14, the scanning optics continually scan across the document 10 in the direction of the arrow 16, which is normal to the document margin 18. If the text margin 20 is parallel to the document margin 18, such a scanning pattern causes each line of text to be scanned along its length as it passes under the optic scanning head. As shown in FIG. 2, however, in the skewed document 10, the text margin 20 is not parallel to the document margin 18 but, rather, is at an angle to it. Thus, the lines of text 12 are at the same angle to the optic scan direction of arrow 16. As a consequence, for example, instead of scanning the first text line 22 and then the second line 24, the right half of line 22 (as FIG. 2 is viewed) is encountered first, followed by the left half of line 22 and the right half of line 24, and finally the left half of line 24. If the data thus scanned are then searched on the normal basis, it yields a meaningless interpretation of the text 12. In a similar fashion, depending on interline spacing and document tilt, the left portion of text line 22 may be encountered on the early part of the scan while the right portion of that line 22 is encountered later in the scan.\nIn certain optical character recognition applications, it is essential that character line rotational effects be removed prior to processing to avoid the generation of meaningless date. Such rotations are usually caused by document feed skewing or typewriter induced skewed text in the actual source document such as that of FIG. 2.\nOne prior art method of circumventing the skew problem is to physically rotate the document or the scanning head in the machine. Such an approach is inconvenient at best, however. What would be preferable, and is therefore the object of the present invention, is a completely non-mechanical correction system and method which identifies and bounds lines of text and which:\n1. Automatically accommodates to arbitrary fonts and sizes thereof.\n2. Automatically accommodates to inter-mixed fonts.\n3. Automatically accepts random inter-character and inter-text line spacings.\n4. Automatically accommodates to a wide range of document skew angles.\n5. Automatically accepts text lines of variable lengths and locations.\n6. Automatically recognizes non-textual data and ignores it in text line finding and bounding."}
{"text": "Wireless communication technologies are used in connection with many applications involving laptop computers, cellular telephones, user equipment, tablets, etc. Wireless communication technologies are tasked with handling increased amounts of data traffic, where the types of data being transported through mobile wireless networks have changed dramatically. This is because of device sophistication, which fosters data-intensive activities such as displaying movies, playing video games, readily attaching photographs to e-mails and text messages, etc. Moreover, video file-sharing and other types of usages (more traditionally associated with wired networks) have been gradually displacing voice as the dominant traffic in mobile wireless networks. This data intensive content burdens the network, as bandwidth is a finite resource. There is a significant challenge for system architects and mobile operators to maintain a stable/reliable network environment and, further, to optimize network resources for engendering acceptable device performance for subscribers."}
{"text": "The recent explosion of the popularity of the World Wide Web (“Web” for short, and hereinafter referred to in the lower case as “web” in the context of an adjective or adverb, e.g., web pages) has made the Internet one of the most important media for mass communication. The Web is used for many applications such as information retrieval, personal communication, and electronic commerce and has been rapidly adopted by a fast growing number of Internet users in a large part of the world.\nUsing the Web, users can access remote information by receiving web pages through the Hypertext Transfer Protocol (HTTP). The information in a web page is described using the Hypertext Markup Language (HTML) and eXtensive Markup Language (XML), and is displayed by software called web browser. Web pages of earlier design are considered static because they do not include any logic that can dynamically change their appearances or provide computations based on user input. Subsequently, the Java™ (Sun Microsystems) programming language was incorporated in web pages in the form of applets. An applet is a small Java™ program that can be sent along with a web page to a user. Java™ applets can perform interactive animation, immediate calculations, or other simple tasks without having to send a user request back to the server, thereby providing the dynamic logic in web pages.\nJava™ is an object-oriented programming language which can be used for creating stand-alone applications. Writing Java™ programs typically requires different and more extensive skills and training than composing web pages. The learning curve for writing Java™ programs is typically longer than that for writing web pages. Not all web page authors therefore are expert Java™ programmers.\nRecently, to make it easier to embed logic in web pages, an easy-to-write script language called JavaScript™ (Sun Microsystems) has been supported by popular web browsers to be incorporated into web pages. JavaScript™, capable of embedding logic for computation based on user input, brings dynamic and powerful capabilities to web pages. JavaScript™, unlike Java™ which is a full-fledged programming language, has a simpler syntax and is much easier to learn. Because of this easy-to-write feature, JavaScript™ has currently become a popular way to embed logic in web pages by many web page authors.\nAlthough JavaScript™ brings easy-to-write logic to web pages, it is limited to browser functions and works with HTML elements only. It can only be used to create simple applications under the contexts of the browser, such as changing the web page's visual presentation dynamically and computing user input quickly without sending a user request back to the server (for such computation). Thus, web pages with JavaScript™ logic cannot be used to create stand-alone applications that require access to a full range of resources on the user's computer such as the file system management and the display area beyond the browser's window. In general, web pages cannot be processed in non-browser contexts.\nAt the present, stand-alone applications are typically written in traditional programming languages (also called 3GL for 3rd Generation Languages) such as C, C++, and Java™, or Fourth Generation Languages (4GL) such as Visual Basic™. Through these languages, stand-alone applications interact directly with operating systems through operating system APIs (application programming interfaces) or indirectly with library functions which may in turn call these operating system APIs. The capability of accessing the operating system APIs gives an application the control of computing resources in a computer.\nIf web pages had embedded logic that could access a whole range of computing resources enabled by these operating system APIs, they could then be used to develop stand-alone applications just like any of the aforementioned 3GL and 4GL languages. Using web pages to develop stand-alone applications would have many advantages. First, web page authors who do not possess the skill and experience in writing 3GL/4GL applications could develop stand-alone applications using the web page technology they profess.\nSecondly, the web technology components that can process the visual presentation language (e.g., HTML), the data modeling language (e.g., XML), and the communication protocol (e.g., HTTP) are available in most computers, which can connect to the Internet through the Web. This would provide an advantage in that using web pages to develop applications, a developer could very efficiently integrate these components. This is because, whereas 3GL/4GL applications can integrate these components programmatically, web pages could integrate them declaratively through languages such as HTML and XML. In general, the shorter learning curve and development time of web pages, as compared with 3GL/4GL programs, would result in a shorter time and lower cost in the development of software applications. The present invention addresses this issue by providing methods and apparatus in a software system that manage the life-cycle of software applications, which are composed of web pages that are not limited to the browser contexts and that have access to the full range of operating system resources.\nAnother issue of the processing of computer software addressed by the present invention is the software installation process. Typically, the installation of a software application is achieved by a special-purpose program which comes with this software and is written only for the purpose of installing this software. This is evident in the existence of a “setup.exe” or “install.exe” program in almost all software packages for PCs (personal computers). This method of software installation means that developers for each software application have to write a specific install program just to install their software.\nIn general, an install program for an application needs to configure a list of settings that are used to establish a proper environment or context for this application before it can be properly installed. These settings may include, for example, the basic operating system setup such as the registry entries, location setup such as the directory or folder in which the application is to be stored, link setup such as the short-cut link to this application, the graphic setup such as the icon of this application, and the dependency setup such as other applications that this application depends on for execution.\nTo properly setup each setting, e.g., one of the aforementioned settings, the install program typically takes the determined value of this setting and processes an action specific to this setting. For example, the registry entry setup action may be to add the determined registry entry values to the proper registry files, whereas the dependency setup action may be to investigate if all applications that the application to be installed depends on are already installed and, if not, to display an error message. Typically, the value of a setting is either determined by user input during the installation process, such as the directory where the application is to be stored, or predetermined by the install program, such as the list of applications that its application depends on.\nIn general, an install program first configures each setting by determining its value (by user input or pre-configuration) and then invokes the setup action for this setting. Because applications may have a different set of pre-configured setting values, each application requires a unique install program. Furthermore, if a new version of an application changes the value of one of its install settings, such as a new icon, the install program for this application has to be rewritten to incorporate this new value.\nIt would be advantageous to the application developers if they did not need to write a new install program for each new version of an application they develop. Instead, it would be desirable, for each version of an application, to construct a list of install settings with pre-configured values for this application using a data modeling language such as XML, which could be provided together with this application for installation. This way, a standardized install program would then be deployed by the user's computer to decode the install settings and values and conduct proper installation for this application based on these values. This standardized install program could then be used to install all applications whose install settings and values are modeled by a language understood by this install program. With many applications installed using a standardized install program, the users would also have a consistent experience in the installation process for all these applications.\nThe present invention addresses this issue by providing methods and apparatus of software installation in which a standardized install manager exists in a computer system to perform the installation process for all software applications whose install settings and values are modeled by a language understood by this install program.\nYet another issue of today's computer software addressed by the present invention is the security management of software applications. Traditional stand-alone applications based on programming languages such as C and C++ typically have access to all the operating system resources through the calling of operating system APIs. In this case, the security context, i.e., the limit of system resource access, for these applications is the entire system. Based on this security context, it is possible that an application can, inadvertently or maliciously, damage not only its own data but those of other applications that share the same computer system.\nIn a virtual machine environment, such as the Java™ Virtual Machine, the security context of an application (such a Java™ program) is defined by the virtual machine. A misbehaving application thus can only create external damage allowable by the virtual machine. However, there can be many different types of applications running on the same virtual machine and while each one of them may have a different security need, they are forced to run under the same security context (that defined by the virtual machine).\nIt would be advantageous if each application had its own security context that is predetermined by the system management policy. Thus, based on its level of security risk, an application could be associated with a security context which regulates the system resources to which this application can or cannot access. This way, a misbehaving program in an application with a restrictive security context would cause minimum damage to the system as a whole. The present invention addresses this issue by providing methods and apparatus of a computer system in which each application has its own security context.\nYet another issue of today's computer software addressed by the present invention is the web cache system for software applications. Web caching is traditionally performed by the web browsers and web proxies whose primary tasks include transmitting web objects over the network. Web pages in the context of a web browser contain hyperlinks to web objects through textual or graphic anchors. The user requests a web object from a web page when this page is displayed by the web browser and the user selects, through the mouse or other pointing mechanism, the anchor of this object.\nWhen a web object is requested through a web browser with the web caching feature, the web browser first checks to see if the object exists in its cache. If so, this object in the browser's cache is returned to the request web page. If the object does not exist in the browser's cache, the browser uses the Uniform Resource Locator (URL) of this object to locate its location in the Internet and retrieves it through a data transfer protocol such as HTTP. When the browser receives this object, it typically displays this object while storing a copy in its cache.\nApplications accessing web objects could be composed using web pages. However, if web pages are processed in the context of the browser, the web objects requested by them in a client computer can only be cached by the browser in the computer. In other word, in a client computer, web page based applications under the browser contexts use only the browser's cache for web caching.\nDifferent web applications however may access web objects with different characteristics. For example, one web application may access web objects that rarely change over time whereas another may access web objects that change highly frequently. It would be advantageous to deploy a sizable space to cache static web objects for the first application while little or no space for the second because any cached objects will be outdated immediately. In general, it would be advantageous that each application has its own web cache.\nFurthermore, traditional web caching by browsers only cache web objects of certain types that are defined in HTTP. Some applications may need to retrieve objects from the Web with types not defined in HTTP. Examples of object types not defined by HTTP may include executable files, spreadsheet files, and documents with proprietary structures. Caching these non-HTTP-defined objects could provide a performance advantage to applications that retrieve objects of these types through the Web.\nThe present invention addresses the issue of web caching for applications by providing methods and apparatus to provide each web application a separate cache for both the HTTP-defined and non-HTTP-defined objects from the Web."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic leakage transformer.\n2. Description of the Related Art\nA magnetic leakage transformer is used, for example, in a discharge lamp operating device used to operate a discharge lamp such as a fluorescent lamp. The magnetic leakage transformer is used in the device not only to generate a high voltage necessary to light a discharge lamp but also to limit a discharge current to be supplied to the discharge lamp as a current-limit inductance while the discharge lamp is operating.\nIf a nonmagnetic leakage transformer is used instead of the magnetic leakage transformer in the discharge lamp operating device, a choke coil to be used as a current-limit inductance must be associated with the nonmagnetic leakage transformer. In this case, the construction of the circuit of the discharge lamp operating device is complicated, the number of components of the discharge lamp operating device is increased so that its assembling is complicated, and the entire discharge lamp operating device is increased in size.\nFIG. 1 schematically shows a conventional magnetic leakage transformer used in a discharge lamp operating device. The conventional magnetic leakage transformer has, as shown in FIG. 2, a core unit 14 in which a pair of cores 10, 12 each having E-shaped plane are combined in a state that extending ends of three legs 10a, 10b, 10c and 12a, 12b, 12c are opposed. A pair of central legs 10a, 102a of a pair of the cores 10, 12 of the core unit 14 are covered with a primary winding spool 16a and a secondary winding spool 18a, and a primary winding 16b and a secondary winding 18b, both of which are consisted of an insulated wire, are respectively wound on the primary and secondary winding spools 16a and 18a. The primary winding spool 16a and the primary winding 16b constitute a primary winding unit 16, and the secondary winding spool 18a and the secondary winding 18b constitute a secondary winding unit 18. The primary and secondary winding spools 16a and 18a respectively have supporting bases 16c and 18c extending along the lower surfaces of a pair of the corresponding cores 10 and 12, and the above-described conventional magnetic leakage transformer is mounted at a predetermined position on a circuit board of a discharge lamp operating device through the bases 16c and 18c.\nIn the above-mentioned conventional magnetic leakage transformer, a pair of side legs 10b and 10c of the three legs 10a, 10b, 10c of one core 10 have the same length, and a pair of side legs of the three legs 12a, 12b, 12c of the other core 12 also have the same length. The central legs 10a and 12a are shorter than the pair of side legs 10b, 10c or 12b, 12c disposed at both sides thereof. As shown in FIG. 1, a pair of the side legs 10b, 10c and 12b, 12c of a pair of the cores 10, 12 are abutted at their extending ends against each other in a state that the pair of cores 10, 12 are associated with each other as described above, and a magnetic leakage gap \"G\" is created between the extending ends of the central legs 10a and 12a.\nWhen a current is supplied to the primary winding 16b of the conventional magnetic leakage transformer constructed as described above, a magnetic flux directed from the central leg 10a of one core 10 corresponding to the primary winding 16b to the central leg 12a of the other core 12 corresponding to the secondary winding 18b is generated in the core unit 14, and this magnetic flux is passed through a magnetic passage returned to the central leg 10a of the one core 10 through the pair of side legs 12b, 12c of the other core 12 and the pair of side legs 10b, 10c of the one core 10. A current having a predetermined relationship to the current supplied to the primary winding 16b is generated in the secondary winding 18b by the magnetic flux. A magnetic resistance generated in the magnetic leakage gap \"G\" between the central legs 10a and 12a of the pair of cores 10 and 12 constitutes a leakage inductance for limiting a discharge current while the discharge lamp connected to the secondary winding 18b is operated.\nThe transmission efficiency of magnetic energy to be transmitted from the primary winding 16b to the secondary winding 18b in the above-described conventional magnetic leakage transformer is determined by the interlinkaging number of exciting magnetic fluxes generated by the primary winding 16b to the secondary winding 18b. Thus, the lesser the magnetic resistance in the magnetic circuit is and the higher the permeability of the core unit 14 is, the higher the transmission efficiency of the magnetic energy becomes and hence the reduction in size of the transformer can be promoted.\nHowever, in the above-mentioned conventional magnetic leakage transformer, the magnetic leakage gap \"G\" increases the magnetic resistance. The magnetic leakage gap \"G\" is necessarily indispensable to prevent magnetic saturation of the above-described conventional magnetic leakage transformer in an inverter operation, but the size of the gap \"G\" required therefor is smaller than that of the gap \"G\" necessary to obtain the leakage inductance.\nIn order to obtain a desired discharge starting voltage for starting the operation of the discharge lamp by compensating a large decrease in the transmission efficiency of magnetic energy generated by the large gap \"G\" necessary to obtain a leakage inductance, in the above described conventional magnetic leakage transformer the numbers of turns of the primary and secondary windings 16b and 18b are larger than those in the above described conventional magnetic leakage transformer. Accordingly, the primary and secondary winding units 16 and 18 are large in size, and hence the entire magnetic leakage transformer is large in size and weight.\nIn order to obtain a desired discharge lamp starting voltage by the above-mentioned conventional magnetic leakage transformer having large magnetic resistance, the value of the exciting current to be supplied to the primary winding 16b must be increased as compared with that of the nonmagnetic leakage transformer."}
{"text": "This invention relates to optically interconnecting opto-electric components on different integrated circuit (IC) chips, and also on the same IC chip, using a printed circuit board (PCB) on which the IC chip or chips are mounted.\nMany electronics systems, including computer motherboards, include one or more IC chips mounted on PCBs. A PCB provides a surface on which the IC chips are mounted, and also provides electrical interconnections between the IC chips.\nSignalling speed requirements between different IC chips in the same electronics system, and perhaps mounted on the same PCB, are ever increasing. In some cases, electrical signaling may not provide the needed, or desired, bandwidth, or may provide the bandwidth with costs. Some of the costs include a more complex design, in terms of multiplexing and demultiplexing the signals into multiple parallel lines. There may also be costs in terms of noise, both because the speed of the signaling may be nearing phyical limits and because of cross-talk between the parallel electrical interconnects.\nOptical signaling, as compared to electrical signaling, offers significantly higher bandwidth and eliminates, or greatly reduces, the noise problems inherent with electrical signaling. An example where bandwidth requirements are making optical signaling between IC chips in the same system increasingly attractive, and in fact may require optical signaling, is in computer motherboards. For example, signaling between processor IC chips and memory IC chips on the same motherboard are already in some systems two gigabytes per second, and will certainly only increase in the future.\nOptical signaling, however, poses design challenges not posed with electrical signaling. For example, optical signaling requires there to be an optical waveguide interconnection between the signal source and detector. In some cases, the optical interconnection between two IC chips within the same system has been provided with conventional optical fibers. However, this approach has its disadvantages. First, the optical fibers add cost to the system. Also, optical fiber connectors are typically large, and thus consume sometimes precious space, and the labor involved in providing connections to optical fibers is typically significant.\nBetter approaches to providing optical interconnects between IC chips within the same system, and even between opto-electronic components on the same IC chip, are therefore needed."}
{"text": "Virtual memory allows programmers to use a larger range of memory for programs and data than is provided by the actual physical memory available to a processor. In addition virtual memory allows programs to be loaded in parallel to one another with a memory map that is impassive to the presence of other programs and the location to which it is loaded. A computing system maps a program's virtual addresses (also known as a Linear Address in the IA32 architecture—in this description ‘linear address’ or LA is used synonymously with ‘virtual address’) to real hardware storage addresses (e.g., physical memory addresses) using address translation hardware. The hardware uses a tree of tables in memory as the input data for the address translation. The root of the tree is pointed to by a register that holds the physical address of the first table in the page table tree. An example of such a register is CR3 in the IA32 architecture. Page table entries (PTEs) are addressed using a base and an index. The base is taken from a register or a previous table memory pointer. The index is using a portion of the linear address. The PTE includes either the page, if the rest of the sub tree is not present in memory, or a memory pointer and other information to be discussed below. The memory pointer is for a page in memory that may either include data (that belongs to the application or the operating system) or another level of the page-table. If it is the later case, another portion of the linear address is used to index into the table in a scheme similar to what is described. If the address is for an application page, the physical address is constructed by adding the remaining bits of the address (that were not used for indexing) to the page base address that was retrieved from the page table entry. Also, some embodiments may instantiate several translation schemes (e.g., different table tree indexing structures) as described in a register, for the sake of simplifying this description we will consider such information part of the CR3 register, even though it may be kept in one or more other registers.\nBeyond the address translation information, these tables include information such as access rights read, write or execute, presence of the data in memory, caching policy indications, page modified state, etc. In some cases, a page table may include pages of different sizes, where larger pages are pointed to in a lower level of the page table tree (instead of pointing to another page of pointers) The size of the page pointed to is stored as an attribute in the page table tree (typically in the level that points to the data page).\nTo retrieve the physical address, the page-table entries are read in a recursive manner starting from the root (CR3 in IA32) and properties of the page are retrieved and merged. The IA32 Programmers' Reference Manual (e.g., Volume 3A) provides an example of a conventional approach to retrieving physical addresses. This process requires several memory access operations and is implemented by Page-Table Handling hardware or uCode sequences.\nOccasionally, software is required to retrieve the physical memory address. In such cases either an emulated full table walk or a shortened heuristic that is based on the limitations of the setup of the table that the operating system imposes are used.\nA translation look aside buffer (TLB) is a cache that holds the result of previous translations such that successive accesses to an address (or a range of addresses) may avoid walking the data structure and can use the results of a previous translation. In many cases the address translation also checks for the operation to meet the conditions set for the memory location. Conventional address translation instructions typically return a physical memory address for a linear address provided as an operand without providing any additional information."}
{"text": "An image display device that includes: a fluorescent screen in which stripe-shaped regions, in which phosphors that are excited by light to emit fluorescent light are formed, are repeatedly formed in the in-plane direction; and a scanning system that scans the fluorescent screen with excitation light is known. Phosphors include a phosphor that emits red fluorescent light, a phosphor that emits green fluorescent light, and a phosphor that emits blue fluorescent light. These color phosphor regions are repeatedly formed in a predetermined order on the fluorescent screen.\nGenerally, in an image display device of this type, the relative positional relationship between the scanning system and fluorescent screen typically undergoes change due to various causes such as vibrations or distortion, changes in the environment such as in temperature or humidity, the effect of gravity, or changes that occur with the passage of time. In addition, if the scanning system is a resonance mirror or the like, the temperature of the resonance mirror that reflects excitation light that has been intensity-modulated varies depending on the intensity of the excitation light. As a result, the resonance frequency of the resonance mirror also varies. If the resonance frequency of the resonance mirror varies, the phase and amplitude of the scanning system vary and thus the scanning position deviates from the correct position.\nIf a static variation and a dynamic variation such as variation of the relative positional relationship between the scanning system and the phosphor screen or a variation of the scanning position occurs, stripe-shaped or matrix-shaped color phosphor regions cannot be irradiated with the excitation light at an appropriate timing. As a result, the luminance of fluorescent light which the individual color phosphor regions emit vary and it causes deterioration of color purity of a displayed image.\nTo solve such a problem, the positions of the individual color phosphor regions on the fluorescent screen need to be accurately detected and the individual color phosphor regions need to be irradiated with excitation light at an appropriate timing.\nPatent Literature 1 describes an image display device that can control the irradiation timing at which phosphor regions are irradiated with excitation light.\nThe image display device described in Patent Literature 1 has a light source, a fluorescent screen, a deflection unit that scans the fluorescent screen with excitation light emitted from the light source, a half mirror located in the traveling direction of the excitation light directed from the deflection device, a photo detector, and a drive circuit that controls the light emission timing at which the light source emits light based on the output signal of the photo detector.\nThe fluorescent screen has individual color (red, green, and blue) stripe-shaped visible fluorescent phosphors that are repeatedly formed at predetermined intervals in the in-plane direction and stripe-shaped dark lanes formed adjacent to individual visible fluorescent phosphors. Stripe-shaped reflection means made of a cube mirror is formed at every second dark lane. The reflection means reflects incident light in the opposite direction of the direction of the incident light.\nThe fluorescent screen is irradiated with the excitation light directed from the deflection unit through the half mirror. The light emission timing is controlled in such a manner that the fluorescent screen is scanned with a predetermined quantity of excitation light. When the fluorescent screen is scanned, the individual color visible fluorescent phosphors and the reflection means are irradiated with the excitation light. The excitation light with which the reflection means is irradiated becomes retro-reflection light that travels in the opposite direction of the direction of the incident light. The retro-reflection light reaches the half mirror. Part of the retro-reflection light reflects on the half mirror and then enters the photo detector.\nThe output signal of the photo detector is supplied to the drive circuit as an index signal that serves to detect the positions of the individual visible fluorescent phosphors. The drive circuit predicts the positions of the individual color visible fluorescent phosphors based on the output signal of the photo detector and controls the light emission timing of the light source such that the visible fluorescent phosphors are irradiated with the excitation light at an appropriate timing."}
{"text": "It has been found, in use of filtering devices of the type disclosed and claimed in my U.S. Pat. No. 3,471,017, that problems can arise in relation to the filtering of polymers having high frictional coefficients and, in particular, polymers which in addition exhibit negligible shrinkage or even positive expansion on solidification, e.g., polystyrene foam (produced for example by the pentane injection process). With such polymers the progressive solidification required to produce efficient sealing plugs in the inlet and outlet ports through which the filter band enters and exits from the filtering passageway would, in the conventional ports (especially the outlet port) as exemplified in my U.S. Pat. No. 3,471,017 cause a progressive setting up of stresses in the continuously forming sealing plugs which, being transmitted to the adjoining walls of the port, might result in the plug binding firmly in the port and, owing to the high frictional coefficient of the material, preventing forwarding movement of the sealing plugs and of the filter band."}
{"text": "Fuel cells which generate electric current by the electrochemical combination of hydrogen and oxygen are well known. In one form of such a fuel cell, an anodic layer and a cathodic layer are separated by an electrolyte formed of a ceramic solid oxide. Such a fuel cell is known in the art as a “solid oxide fuel cell” (SOFC). Hydrogen, either pure or reformed from hydrocarbons, is flowed along the outer surface of the anode and diffuses into the anode. Oxygen, typically from air, is flowed along the outer surface of the cathode and diffuses into the cathode. Each O2 molecule is split and reduced to two O−2 anions catalytically by the cathode. The oxygen anions transport through the electrolyte and combine at the anode/electrolyte interface with four hydrogen ions to form two molecules of water. The anode and the cathode are connected externally through a load to complete the circuit whereby four electrons are transferred from the anode to the cathode. When hydrogen is derived from “reformed” hydrocarbons, the “reformate” gas includes CO which is converted to CO2 at the anode via an oxidation process similar to that performed on the hydrogen. Reformed gasoline is a commonly used fuel in automotive fuel cell applications.\nA single cell is capable of generating a relatively small voltage and wattage, typically between about 0.5 volt and about 1.0 volt, depending upon load, and less than about 2 watts per cm2 of cell surface. Therefore, in practice it is usual to stack together, in electrical series, a plurality of cells. Because each anode and cathode must have a free space for passage of gas over its surface, the cells are separated by perimeter spacers which are vented to permit flow of gas to the anodes and cathodes as desired but which form seals on their axial surfaces to prevent gas leakage from the sides of the it stack. The perimeter spacers include dielectric layers to insulate the interconnects from each other. Adjacent cells are connected electrically by “interconnect” elements in the stack, the outer surfaces of the anodes and cathodes being electrically connected to their respective interconnects by electrical contacts disposed within the gas-flow space, typically by a metallic foam which is readily gas-permeable or by conductive filaments. The outermost, or end, interconnects of the stack define electric terminals, or “current collectors,” which may be connected across a load.\nA complete SOFC system typically includes auxiliary subsystems for, among other requirements, generating fuel by reforming hydrocarbons; tempering the reformate fuel and air entering the stack; providing air to the hydrocarbon reformer; providing air to the cathodes for reaction with hydrogen in the fuel cell stack; providing air for cooling the fuel cell stack; providing combustion air to an afterburner for unspent fuel exiting the stack; and providing cooling air to the afterburner and the stack. A complete SOFC assembly also includes appropriate piping and valving, as well as a programmable electronic control unit (ECU) for managing the activities of the subsystems simultaneously.\nIn an SOFC being supplied with fuel from a reformer, the fuel cell supply gas is provided directly from the reformer. The reforming process takes place at an elevated temperature (800° C.-1000° C.) that is somewhat higher than the optimum stack operating temperature. For proper operation of the stack, it is preferable that the anode gas be at a temperature somewhat below the stack operating temperature, preferably between about 550° C. and 700° C. In addition, it is preferable that the inlet temperature of cathode air be about the temperature of the anode gas.\nIt is a principal object of the present invention to provide optimal tempering of anode gas and cathode air."}
{"text": "The present invention relates to an emergency-stop circuit, which is an integral part of the typical industrial machine. More particularly, this invention relates to a centralized switching system and method for an emergency stop circuit.\nIn industrial equipment, the traditional emergency-stop circuit consists of a xe2x80x9cself-latchingxe2x80x9d relay that contains a number of closed (kill) switches which are connected in series, and when any one of the switches is opened, the relay is de-energized. Power is restored when all kill switches are closed, and a xe2x80x9cmotors-onxe2x80x9d momentary switch (e.g., push-button switch) manually closes the contacts of the relay. The relay contacts are the last link in the serial chain of switches that energizes the coil of the relay. It is self-latching in the sense that when the motors-on switch is released, the contacts are in the coil energizing circuit that keep them closed in the first place. The coil energizing circuit is referred to herein as the emergency-stop circuit.\nA robust, traditional circuit may have many kill switches in the emergency-stop circuit. These switches are typically distributed all over the machine. For example, lever-type switches are installed on door panels, so that power is killed (i.e., shut off) when one of the doors opens. This is referred to as the normally open configuration (NO), which means that the switch must be tripped to conduct. This kind of kill switch is the first to be defeated in practice. It is often taped or strapped closed so that a door may remain open during operation of the machine. (A common purpose for the defeat is debugging by a maintenance technician.) When there are several doors defeated in this manner located throughout a large machine, the probability is higher than desirable for a maintenance technician to inadvertently leave a switch defeated and return the machine to what will be unsafe use. Also, the cycle of taping/strapping and removal thereof causes wear and tear on the lever-type switch for which it was not designed.\nOther types of kill switches used in the industry include over-travel switches. These switches normally operate in the closed configuration (NC), which means that tripping of the switch opens the circuit. These switches include lever-type, magnetic, infrared, or the like. To defeat over-travel switches, the switches are temporarily removed, terminals jumpered, mounting screws loosened, and brackets are slid out of the way. This also creates opportunity for mistakenly leaving kill switches defeated (or misaligned) throughout the machine when it is returned to service.\nAnother example of a kill switch is an air pressure switch sensing an air line that delivers required air to an air bearing spindle. In a demonstrating test, or debug mode, the machine may be run without the spindle running (no air supplied or air temporarily unavailable). This requires the jumpering of the kill switch during such time. Afterwards, forgetting to re-enable the switch allows running of the spindle without air, which leads to hardware damage.\nEvidently, safe use of the traditional emergency-stop circuit requires experience and diligence on the part of the maintenance technician who attempts to temporarily bypass sections of the circuit in order to test or debug the system. Oversight due to distribution of the switches over numerous parts of the machine/device can cause him to forget to re-enable a kill switch before returning equipment back to duty.\nAdditionally, in order to test and debug, the technician must also disable certain devices whose power is controlled by the emergency-stop circuit. There is no straightforward, universal way to do this other than disconnecting the power to the device. This may be easy in some cases or not possible, very cumbersome, or unsafe in others.\nA final consideration for these testing and debugging methods is the time required for a technician to trace through a machine in order to determine where to disable a kill switch or where to disconnect power to a device. Additionally, managerial time may be spent generating documentation in order to aid the technician\"\"s task. This becomes apparent when one considers a factory floor that possesses a vast array of one-of-a-kind machines, all of which utilize some variant of the traditional emergency-stop circuit. Here, hypothetically, each circuit possesses essentially the same topology but utilizes different components that are located in different places and connected by a slightly different wiring scheme.\nIn spite of this, implementation of traditional emergency-stop circuits that are intrinsically xe2x80x9csafexe2x80x9d is certainly feasible and has been done for many years. There are reasons for the apparent success. It is a simple circuit, even though it is distributed throughout the machine. It well established. There are few components. But these are also the reasons why the circuit has not matured.\nTypically, experienced engineers are reluctant to add new parts and kill switches to the circuit in an effort to xe2x80x9ckeep it simple.xe2x80x9d In developing prototypes or one-of-a-kind machines, important kill switches such as a watchdog circuit and a computer ready are often omitted. Also, some kill switches having solid state outputs (e.g. NPN) do not fit into the serially connected topology. Each requires an extra part, such as an intermediate electro-mechanical relay, whose contacts are in the kill switch chain, and whose coil is controlled by the solid state output. Because of this, sensors employing solid state outputs are avoided, and their less reliable mechanical counterparts are used instead.\nEssentially, there is a mindset among skilled engineers concerning the altering of the traditional circuit\"\"s topology. Typically, the skilled engineer begins a new project assuming that he will use the traditional circuit. Valuable time is spent on other areas and is not devoted to re-engineering the architecture for the traditional circuit or evaluating its expanded role in the project. In fact, it is not obvious to the skilled engineer to change the traditional circuit in any way in order to add functionality that can be safely incorporated within it. Such functionality, if implemented, is therefore left to be distributed throughout the remainder of the system, intermingled with unsafe subsystems such as the computer.\nWhen implemented, for example, secondary outputs, such as amplifier xe2x80x9cenablexe2x80x9d or xe2x80x9cinhibitxe2x80x9d signals, are not usually incorporated into an emergency-stop circuit. If driven at all, a software program running on a computer having optically isolated digital outputs usually drives them. Furthermore, other feedback signals, such as xe2x80x9cstatusxe2x80x9d or xe2x80x9cfaultxe2x80x9d signals, are not used in emergency-stop circuits as kill inputs. This is generally because each signal is in a non-conducting state when the circuit is killed, which prevents the traditional circuit from restarting. If used at all, these feedback signals are likewise connected to the computer for the purposes of monitoring.\nDesigning in this way fosters subtle system-wide shortcomings, which can permit potentially unsafe or undesirable operation. Resulting failures or odd performance is not attributed to the emergency-stop circuit, since its simple circuitry and lack of substantial functionality are not directly responsible. Consequently, effort is typically not expended to evaluate its functionality.\nOne of the shortcomings becomes apparent when the traditional system enters into a power-loss period, which generally begins when the emergency-stop circuit is killed and ends when all residual power has been dissipated. During this brief period (e.g., 2 sec.), uncontrolled motion of motors can occur for some designs, because the motors are not being controlled, yet they are still technically powered by residual power in the system. In order to suppress this, designers have used the computer-controlled secondary outputs (enable, inhibit) in conjunction with the emergency-stop circuit to simultaneously cut power and disable the connected devices. This works in most cases, but is tedious to design, not flexible, and application specific. One case when this design fails is when the building power fails, which causes the computer to also cease functioning. Here the inhibit signal may not get to the device, which again creates an environment for briefly uncontrolled motion.\nMost of the examples found in existing technology are concerned with passive monitoring of the emergency-stop circuit. This approach is useful in determining which kill input was responsible for stopping the circuit, but it does not provide any configuration options for startup or power-loss periods. The following patents, each of which is incorporated herein by reference, demonstrate this approach: U.S. Pat. No. 4,263,647 to Merrell, et al, entitled xe2x80x9cFault Monitor for Numerical Control Systemxe2x80x9d; U.S. Pat. No. 5,451,879 to Moore, entitled xe2x80x9cElectromechanical Relay Monitoring System with Status Clockingxe2x80x9d; U.S. Pat. No. 4,616,216 to Meirow, et al., entitled xe2x80x9cEmergency Stop Monitorxe2x80x9d; and U.S. Pat. No. 5,263,570 to Stonemark, entitled xe2x80x9cConveyor Belt Emergency Stop Indicator Light System.xe2x80x9d Configuration options do exist in the above noted patents but only in the form of providing cascaded inputs and outputs so that multiple groups of sensors may be monitored. Other patents of interest include the following: U.S. Pat. No. 4,912,384 to Kinoshita, et al., entitled xe2x80x9cEmergency Stop Control Circuitxe2x80x9d discloses the traditional active portion of the emergency-stop circuit; U.S. Pat. No. 5,319,306 to Schuyler entitled xe2x80x9cPortable Electrical Line Tester Using Audible Tones to Indicate Voltagexe2x80x9d discloses circuits that provide audio status in the form of line testers, where the leads are brought into contact after the line is energized to check it.\nTraditional approaches to supplying power to motors during a power-loss period (period beginning with the loss of AC motor power and ending with either the total loss of all stored DC motor power or the loss of regulation of any associated logic power supply, whichever comes first) have focused on coarse (non-servo) control or decelerating motors to full stop. However, no approach exists that relates to fields employing emergency-stop circuitry.\nOther patents in this general field are also noted. For example, U.S. Pat. No. 5,278,454 to Strauss, et al. discloses an invention related to the heating, ventilation, and air conditioning field. It describes a motion control system that senses a loss of incoming power and utilizes a dedicated pre-charged circuit to act as a short duration power supply to effect gross motion of a motor to close a damper. U.S. Pat. No. 5,426,355 to Zweighaft, et al., entitled xe2x80x9cPower-Off Motor Deceleration Control Systemxe2x80x9d discloses an invention related to the tape drive industry in which a motion control system whose amplifier stores a dedicated internal PWM signal responsible for supplying open-loop deceleration commands for a given configuration of the tape drive system that is experiencing a power-loss period. U.S. Pat. No. 4,481,449 to Roda entitled xe2x80x9cPower Fail Servo Systemxe2x80x9d discloses an invention that also relates to the tape drive field which describes the use of several xe2x80x9cpower failxe2x80x9d signals that work in harmony to decelerate the motor towards full stop and uses the technique of dynamic braking to harness excess power in the storage capacitor. A signal exists in this example which monitors the logic power supply and appropriately disables (free wheels) the motor once the supply is out of regulation.\nThe present invention solves the problems in the art by providing a centralized programmable emergency-stop circuit that controls the flow of the power necessary for a machine to move its working elements. The invention possesses various levels of programmability that facilitate use of the same circuit across a wide variety of industrial applications and designs, as well as across a wide variety of operational scenarios for the same machine.\nThe circuit of the present invention includes various types of custom programmable kill inputs. These inputs are signals that, subject to their programming, can kill an energized emergency-stop circuit or prevent a killed circuit from energizing (startup). A given kill input can also be programmed to be ignored totally, to kill when inactive, or to also prevent startup when inactive. A given kill input can be programmed so that it only affects the energized circuit and does not restrict startup, and consequently, it may be inactive at startup. Such a programmed kill input is referred to herein as a xe2x80x9cfalling-type,xe2x80x9d because once it does go active, it is the active-to-inactive or falling transition that kills the circuit. Additional programming for the kill inputs exists such as digital filter parameters, clock selection, and the like, as well as time-out options for the falling-type kill inputs, which require them to go active within some period after startup.\nThe present invention also provides programming options to specify conditions for a motors-on signal to energize the circuit and for the control of secondary outputs. While the primary output of the circuit controls the flow of bulk power to working elements, it is the secondary outputs that connect in parallel to the working elements in order to inhibit or enable them. The method of programming secondary outputs determines their behavior, i.e., whether they are disabled entirely for the session, enabled only when the circuit is energized, or enabled based on one of the kill input signals. This latter setting permits a computer to keep a device enabled during a power-loss period, so that a reactionary movement can be effected which drains residual power left in the dying system.\nIn order to improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves performance during the period immediately following the application of electrical power needed to power circuit logic. Specifically, it is the object of the invention to inhibit energizing the circuit for a prescribed interval of time. Additionally, it is the object of the invention to provide programmability so that the interval may be changed.\nIn order to further improve performance during the period immediately following the application of electrical power needed to power circuit logic, it is the object of the invention to provide additional features and programmability. Specifically, it is the object of the invention to provide circuitry that determines whether the circuit has been energized at least once. Furthermore, it is the object of the invention to provide further additional circuitry that drives a dedicated power-up/reset error code which indicates electrical power has just been applied to the circuit logic. The power-up/reset error code therefore supersedes the conventional error code that is generated from all possible kill input sources. Additionally, it is the object of the invention to provide a clear signal capable of clearing the power-up/reset error code (so that the conventional error code may be revealed) and also capable of refreshing conventional error codes thereafter. It is also the object of the invention to provide programmability so that a set of clear input sources may be pre-selected from all available input sources.\nFinally, in order to further improve performance during the period immediately following the application of electrical power needed to power circuit logic, it is the object of the invention to provide additional features and programmability. Specifically, it is the object of the invention to employ a start signal that when inactive inhibits the initial energizing of the circuit. Activation of the start signal occurs in response to the final cycle of a specified number of deactivation and reactivation cycles of a ready-type input signal, and deactivation of the start signal occurs when the circuit is energized. Additionally, it is the object of the invention to provide programmability so that (1) the ability of the start signal to inhibit energizing is optional, (2) the specified number of cycles can be adjusted, and (3) a set of ready-type input signals may be pre-selected from all available input sources.\nIt is also the object of the invention to further employ the same start signal in subsequent energizing cycles in order to further improve performance. Specifically, a second specified number of deactivation and reactivation cycles is required in order to activate the start signal. Additionally, it is the object of the invention to provide programmability so that the second specified number of cycles can be adjusted.\nIn order to improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves how the circuit is commanded to energize. Specifically, it is the object of the invention to provide for additional nominal requirements for the activation of a motors-on signal, such as (1) requiring it to be previously inactive and (2) requiring it to be active for a prescribed interval or longer. Additionally, it is the object of the invention to provide programmability so that (1) the interval may be changed, (2) the requirement to be previously inactive is optional, and (3) a set of motors-on-type input sources may be pre-selected from all available input sources. Finally, it is the object of the invention to provide programmability so that (1) a set of monitor contact-type input sources may be pre-selected from all available input sources, where each monitor contact signal is active when the circuit is killed and the associated, downstream monitored relay has fully disengaged and (2) the requirement for a given monitor contact signal to be active for the motors-on signal to be active is optional.\nIn order to further improve the manner in which the circuit is energized, it is the object of the invention to employ a second start signal that when inactive inhibits the energizing of the circuit. Activation of the start signal occurs when all kill input sources are active, where programmability provides for a set of kill sources to be selected from all available input sources. Deactivation of the start signal occurs when the circuit is energized or when one or more of the kill input sources become inactive. Additionally, it is the object of the invention to provide status for the start signal. Furthermore, it is the object of the invention to accommodate watchdog-type kill input sources that toggle on-and-off repeatedly at a rate faster than a prescribed value, where the toggling is the requirement for the watchdog-type kill input to be active. It is also the object of the invention to provide programmability for this so that (1) the requirement for toggling is optional and (2) the minimum rate is programmable. Finally, it is the object of the invention to include in the generation of the start signal an additional, dedicated kill input source that indicates whether an internal circuit error exists.\nIn order to further improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves performance during the period immediately following energizing (right after it is started). Specifically, it is the object of the invention to provide audio status for a prescribed interval. Additionally, it is the object of the invention to provide programmability so that the interval may be changed.\nIn order to improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves the manner in which the circuit is de-energized (killed) or prevented from energizing. Specifically, it is the object of the invention to employ a kill signal that when active de-energizes the circuit or prevents it from energizing. Activation of the kill signal occurs when one or more kill sources become inactive, where programmability provides for a second set of kill sources to be selected from all available input sources. Deactivation of the kill signal occurs when all kill sources from the second set become active. Additionally, it is the object of the invention to include in the generation of the kill signal an additional, dedicated kill input source that indicates whether an internal circuit error exists.\nIn order to further improve performance for the manner in which the circuit is de-energized (killed) or prevented from energizing, it is the object of the invention to provide additional programmability so that pre-selected additional input sources can be dynamically added to the second set of kill sources at some point of time after the circuit becomes energized and subsequently removed at such time that the circuit is de-energized. A given, dynamically added input source may be programmed to be added immediately after the input source becomes active. Additionally, or alternatively, it can be added after a prescribed interval of time following the energizing of the circuit. It is also the object to provide programmability so that this prescribed interval can be adjusted.\nIn order to further improve performance for the manner in which the circuit is de-energized (killed) or prevented from energizing, it is the object of the invention to provided additional programmability so that one of the dynamically added input sources is dedicated to sensing the presence of the bulk power controlled by the circuit. Additionally, it is the object that this input source is an alternating-current type that generates a strobing signal indicative of the active state of the bulk power, where the strobing occurring at a rate faster than a prescribed value is the requirement that the kill input source is active. Finally, it is the object that the minimum rate is programmable.\nIn order to further improve an emergency-stop circuit that controls the flow of bulk power needed for a machine to move its elements, it is the object of this invention to provide additional features and programmability that improves performance during the period immediately following de-energizing (right after it is killed). Specifically, it is the object of the invention to inhibit the re-energizing of the circuit for a prescribed interval of time after it is killed. Additionally, it is the object of the invention to provide programmability so that the interval for the dying period may be changed. Also, it is the object to provide audio or visual status during the dying period.\nIn order to further improve an emergency-stop circuit whose primary output controls the flow of bulk power needed for a machine to move its elements and whose secondary output controls the enable or inhibit of an element, it is the object of this invention to provide additional features and programmability for the circuit so that the source of the secondary output may be selected from a set of available sources. Specifically, it can be selected from the following sources: (1) none so that the element is always disabled, (2) from a signal that is active when the circuit is energized so that the element is enabled only when the circuit is energized, or (3) a dedicated enable-type input source, so that the element is enabled whenever the enable-type input source is active. It is also the object of the invention to provide additional programmability for the third case, which places a programmable pair of restrictions on when the enable-type input source has an effect so that it is used when (1) the circuit is energized or in the dying period that immediately follows de-energizing and otherwise, the element is disabled and (2) a watchdog-type input source is active and otherwise, the element is disabled. The requirement for the watchdog-type input source to be active is that it must toggle on-and-off repeatedly at a rate faster than a prescribed value. Finally, it is the object of the invention to provide additional programmability so that (1) the minimum rate for the watchdog-type input is programmable, (2) the enable-type input source may be pre-selected from all available input sources, and (3) the watchdog-type input source may be pre-selected from all available input sources.\nAccordingly, it is the object of the present invention to provide a programmable emergency-stop circuit that allows various options for the manner in which kill inputs affect the system and further provides options for the manner in which outputs are activated and deactivated. Furthermore, it is an object of the invention to provide programmability to specify the manner and timing for dynamically adding a given input source to the active set of kill inputs. Finally, it is an object of the invention to emp e circuitry that generally avoids software or a microprocessor, so that new functionality coupled with programmability may be safely incorporated within the emergency-stop circuit.\nOne important feature of the invention is its state machine, which provides a framework from which the invention operates. Defined by a set of internal signals that includes start and kill-type signals, the state machine specifies when the circuit may be energized, when it is killed, and when startup is inhibited. The internal signals are generated as a programmable function of time and input source states. Other features include audio status for startup and kill, requirements for startup that ensures desired energizing, requirements for a computer ready signal that ensures synchronization with software running on a computer, provisions for a dedicated error-code that identifies power glitches, and the safe oversight of a power-loss period during which a servo-controlled reflex action may be implemented.\nThe primary advantage for using the invention is that a centralized single circuit can be programmed and employed in a wide variety of machine designs. For a given machine design, for example, the circuit can be reprogrammed and thereby adapted to a different set of operational scenarios. When designing a machine or a plurality of machine/devices, the designer is able to associate any given input source with a desired kill input type that specifies how the input source affects the system. Furthermore, once operational in the field, for example, the machine will require maintenance, and to assist this, the circuit can be definitively reprogrammed from a central location so that certain inputs are temporarily but safely ignored and certain outputs are forced disabled during the maintenance operation.\nOther advantages of the invention are related to timing, filtering, and synchronization. One such advantage is the accuracy, and hence repeatability, that can be applied to timing the motors-on button\"\"s active period as well as to the timing of the start-up delay that prevents the immediate re-start during the DYING state of a freshly killed circuit. The use of timing and other related digital filters significantly reduces the susceptibility of the circuit to background noise. It is also an advantage from a system performance standpoint that the emergency-stop circuit causes the computer program and, thereby, the entire system to be in synchronization via several novel methods.\nThe invention will now be described, by way of example and not by way of limitation, with reference to the accompanying sheets of drawings and other objects, features and advantages of the invention will be apparent from this detailed disclosure and from the appended claims. All patents, patent applications, provisional applications, and publications referred to or cited herein, or from which a claim for benefit of priority has been made, are incorporated by reference in their entirety to the extent they are not inconsistent with the explicit teachings of this specification."}
{"text": "This application is the national phase under 35 U.S.C. xc2xa7371 of PCT International Application No. PCT/US98/24404 which has an International filing date of Nov. 17, 1998, which designated the United States of America.\n1. Field of the Invention\nA method and composition of matter for use as polymeric topcoats for articles and vehicles, such as, aircrafts, naval vessels, clothing and other industrial applications. With regard to an aircraft, xe2x80x9ccold-soakxe2x80x9d of the aircraft wing fuel tank leads to localized wing ice formation under certain environmental conditions. Also, ice forms on the xe2x80x9cpleading edgesxe2x80x9d of the aircraft which detach and enter the jet engines or otherwise influence aerodynamic performance of aircraft wings. Conventional polymer paints and coatings contain a volatile organic content (VOC) that is under increasing regulation by EPA.\n2. Background of the Invention\nSince 1986, the limit for volatile organic content (VOC) of aerospace topcoatings as set by California Rule 1124 has dropped from around 700 g/l to its present limit of 420 g/l or even lower values. Increasing concern over the impact of organic compounds on the quality of life and environment can be expected to lead to further reduction in permissible VOC in coming years. Achieving durable, functional coatings that comply with the VOC regulations and satisfy functional coating requirements is becoming challenging for the aircraft industry and coating suppliers.\nThe southern California environmental control agencies require a maximum of 420 grams/liter of volatile organic compounds (VOC) from a coating material. The cyclic prepolymer coatings will reduce the VOC emissions during coating operation to less than 1 gram/liter of coating material. These new coating processes will provide a coatings technology that is environmentally compliant for the future, whereas existing solvent-borne technologies are compliant on a year-to-basis with a questionable future.\nConventional aircraft coatings used on commercial and military aircraft can be either water based or solvent based. Solvent based coatings are the most widely used. Typical solvents such as xylene, toluene and chlorinated aliphatic hydrocarbons, are required in order to control drying times, pigment distribution and surface smoothness of these coatings. These compounds all have unacceptably high VOC. Furthermore, xylene is a carcinogenic compound and the others are suspected hazardous materials both of which present serious employer liability issues. The EPA is strongly advocating a reduction of all solvents with the exception of water to reduce VOC and eliminate potential carcinogens. A water based coating is a natural alternative and has been developed for primer coatings but has yet to produce satisfactory performance as a topcoating. They contain small but significant amounts of VOC.\nHistorically, coating formulations meet the requirements by using xe2x80x9cexemptxe2x80x9d solvents, or by reclassifying coatings into other categories. Newer approaches for formulations and applications of coatings represented by the approach of the present invention can achieve a reduction of VOC well below 420 g/l, perhaps approaching as low as 0 g/l. This is achieved by using a new polymeric coating technology that will meet the most severe restrictions that are anticipated in the year 2010 less than 100 g/l.\nHigh solids deposition processes are based on water reducible, flame and plasma spray coating processes to implement low VOC coatings and deposition processes through a highly focussed research and development program.\nPlasma spray and flame spray processes and flourinated polymer coatings have advantages because current solvent systems have definite limits for reduction of VOC. Although water-reducible systems have potential for further VOC reduction, they have a significant VOC content and may also exhibit adherence problems. A xe2x80x9csuper-critical fluid spray coating systemxe2x80x9d is capable of reducing VOC by 30-70 percent depending on the type of resins and polymers in the parent coating system. However, the equipment is expensive, complex and bulky, and the pigmentation of coatings using this process is limited.\nPlasma spraying consists of depositing a coating by flowing a powder coating-inert gas mixture through an electric arc plasma. The thermoplastic powder liquefies and flows on the surface. The advantage of this coating process over air spraying of solvent-borne and water-borne coatings is that no solvent or VOC is produced. Also, many materials can be applied with low surface energies, such as chloro- and fluoropolymers which cannot be air sprayed. The disadvantage is that the process produces an ignition source which is hazardous around aircraft and flammable vapors and liquids. The actual cost of the plasma spray coating process is higher than conventional coating processes, but the service life of the plasma sprayed coating is longer and the coating can be thicker to compensate for wear. Lifecycle costs may be lower than conventional coatings. Typical foot-wear on the surface will not damage these coatings.\nThe plasma spray process is a mature technology and equipment is available for use. These coatings can be applied directly to aircraft aluminum surfaces to provide a non-icing surface. Limited colors are available in stock powders, but can be formulated for any color. In order to achieve an optimal coating it also will be necessary to formulate binders and pigments with specific properties.\nFor a thixotropic powder, particles need to coalesce quickly, (tc needs to be short) since xcex7 is time dependent after deformation. The instantaneous viscosity, xcex7, and particle radius, rp, should be small and the surface tension, xcex3, large. For flattening, xcex7, and particularly rp should be small; and xcex3 and particularly h should be large. (Additives may be able to reduce the surface tension, but viscosity seems the primary driver.) Low xcex7 necessitates low molecular weight and higher temperatures, or slower catalysis rate.\nSpecific flatteners, pigments and other additives are necessary to make an effective topcoat from a resin (binder) promoting coating adhesion, providing ultra-violet (UV) radiation protection and color. These additives must be balanced against the requirements for coalescence and flattening, as increasing content of particulate in the coating increases viscosity. Several specific texts on paint chemistry for the production of a topcoat are available to guide coating formulation.\nTwo generic types of coatings are relevant, aircraft topcoatings and industrial maintenance (IM) coatings. Requirements for aircraft topcoatings are stringent. Typically, aircraft topcoat requirements are specified by the military specifications MIL-C-83286B, xe2x80x9cAliphatic Isocyanate Urethane Coating for Aerospace Applicationsxe2x80x9d, MIL-C-85285, xe2x80x9cHigh Solids Polyurethanesxe2x80x9d, or Boeing Military Specifications such as BMS 10-60, xe2x80x9cProtective Enamel.xe2x80x9d\nAn EPA reports summarizes the competitive low VOC coating processes and chemistries available in 1991; the principal ones being powder, waterborne, radiation curable and high solids coatings. The summary of this older reference still appears to represent a good economic and technical assessment of coating possibilities. This report also emphasizes that VOC from coating stripping operation is also considered one of the VOC consequences of the selection of method of coating. Table 1 summarizes coating/application methods and issues in this report.\nBesides the genre of application for coatings, the actual deposition method is an important element in controlling VOC. High solids coatings, for example, achieve low VOC by eliminating the solvent classically used for coalescence, flow and flattening. They rely instead on mechanisms such as thermal or kinetic energy to achieve these ends water-based coatings replace organic solvents with water.\nThe EPA sets forth VOC requirements for industrial maintenance coatings, namely, primers, sealers, topcoats, etc., used in outdoor aggressive environments on structures such as bridges, ships, and hydraulic structures. The proposed VOC limit was 350 g/l; and the 2004 limit (proposed) was 300 g/l.\nThis reports describes in detail the VOC measurement methodology (EPA Reference Method 24, a distillation of ASTM standard test methods) and describes the calculation of VOC emissions. Manufacturers claim that the ASTM D-2369 can produce inordinately high VOC levels, particularly in marine and architectural coatings, as it requires the coating to be heated to 110 C. (230 F.) where excessive loss of volatile components by coating decomposition may occur.\nCamouflage topcoatings must meet low VOC requirements and also very stringent chemical agent resistance requirements. These requirements consist of resistance to chemical decontamination/wash as well as other severe requirements.\nThe baseline coating is a two components solvent-borne polyester/polyisocyanate binder system that is lead, chromium, 1,1,1 trichlorethane free. This candidate new xe2x80x9clow VOCxe2x80x9d coating is a waterborne/dispersible/reducible coating using polyisocyanates and polyesters from Miles, Inc. including Bayhydrol XP-7044 WD polyester, Bayhydur XP-7007 WD polyisocyanate and de-ionized water reducer as needed.\nTypical fillers are cobalt green spinet, chromium oxide, magnesium ferrite and carbazole violet pigments with diatomaceous silica, magnesium silicate and amorphous silica extender pigments. These pigments are added to the polyester component and polyisocyanate is diluted with suitable solvent to meet viscosity of both components and meeting stoichiometry.\nThis coating met all requirements of specification except CAR (chemical agent resistance). VOC is estimated to be xcx9c300 g/l. A fundamental problem of these WB/WD/WR coatings is film porosity that allows chemical agents to penetrate the coating. The CPVC (critical pore volume content) appears to determine gloss. A high CPVC value in the candidate coating is a problem for CAR. The author of the study cites the strategy for improving the CPVC is use of additives to improve wetting/flow/dispersion in this WB/WD/WR system.\nIt would appear that a well designed, low VOC aircraft topcoat may also met CARC requirements.\nTable 2 outlines the basic performance characteristics of aircraft coatings.\nRegulatory bodies tend to restrict the use of deposition systems that have low transfer efficiency. The California AQMD regulations require minimum transfer efficiencies of 60-85% and maximum gun tip gas pressure of 10 psi. Currently, only HVLP and electrostatic spray processes can meet these requirements.\nExisting levels of corrosion protection should be maintained with new low VOC topcoatings. Corrosion (aqueous) requires the presence of water, cations and oxygen. strontium chromate is an important additive to inhibit corrosion. Coating strategy has been to achieve a physical barrier between the substrate and the external environment to prevent moisture and radiation induced coating degradation.\nSince moisture egress is a virtual certainty, coating adhesion becomes a very important coating characteristic. The mechanism of adhesion is either chemical or physical. Although chemical pretreatment of the surface substrate can enhance secondary chemical bonding, and in some cases even achieve primary bonding, the major adhesion mechanism is the mechanical interlocking of the coating with the microscopic surface roughness created in anodizing.\nThe practical lifetime of a military or commercial coating is 4-8 years. This lifetime requirement imposes significant demand for resistance to environmental degradation.\nTraditionally an epoxy primer and polyurethane topcoat are used for aircraft applications. Epoxy primer/polyurethane topcoatings are highly refined to meet the military requirements. Since epoxides are brittle and have very low UV stability, they are used as a primer and the external coating provides the UV protection. The epoxide coatings provide superior resistance to moisture penetration and subsequent corrosion. The combination of the two also has very low water absorption, vapor transmission rate and UV resistance.\nAliphatic isocyanate and polyester are highly developed UV resistant topcoatings and their literature is well documented. Typical aircraft topcoats have a dry film thickness of 0.002+/xe2x88x920.0003 inch (50.8+/xe2x88x927.8 micrometers). Set and hard dry time is typically 2 and 6 hours, respectively. Fully developed properties may not be attained until about 7 days aging.\nFillers such as talc and mica are used to provide an oriented distribution to serve as secondary radiation and physical barriers, silica and metal silicates, carbonates and sulfates are added as physical fillers that reduce gloss and increase opacity.\nWater-borne coatings are one approach to compliant coatings. The primary strategy is to achieve a solution of emulsion of polymer powders whose surfaces are modified with hydrophilic groups. The major difficulty with water-borne coatings is that the use of water as a solvent leads to more porous coatings and adhesion problems related to organic surface contamination.\nHigh solids or powder coatings are also one route to achieving low VOC levels. One approach is to reduce the solvent content of the coating, but this shortens pot life and greatly increases viscosity. These factors increase surface roughness. By moving to lower molecular weight resins, one can achieve improved viscosity and flatter coatings. However, polyisocyanate cured powders have shorter pot life and reduced flexibility due to the more rapid and extensive cross linking due to the lower molecular weight. One strategy to improve this is to use polymers with very narrow molecular weight distribution.\nAn EPA study and a follow-up publication evaluated six (6) coatings. The six types are solvent-borne polyurethane, waterborne epoxy primer w/latex topcoat, solvent alkyd primer/waterborne acrylic, 2 component polysiloxane topcoat, water reducible alkyd primer/acrylic topcoat and solvent alkyd primer/solvent alkyd enamel (standard baseline). The study compared impact, adhesion, pencil hardness and solvent tests and outdoor exposure tests of these coatings. The VOC data from candidate coatings in this study is useful for aerospace topcoats. Of particular interest in this assessment of IM coatings is the determination that a two component polysiloxane coating with a low VOC of about 84 g/l performed extremely well. These two studies showed the polysiloxane coatings exhibited the best VOC levels and performance in environmental testing.\nA second study is also grouped in the IM coatings discussion. Although the study intended to coat F-15 aircraft, only ground vehicles were coated. It consisted of an evaluation of supercritical spray coating and a high pressure-low volume (HPLV) process called ULV (ultra low volume) spraying. The polyurethane coatings had a baseline VOC of about 420 g/l which is too high for current requirements. The study found the supercritical coating process to be unacceptable for field use, and found that the ULV process reduced emissions by about 50%, primarily by reducing the total paint sprayed in the coating process. This result shows that both the coating process itself as well as the formulation of the coating can have significant impact on total VOC emission. This study attempted to spray high solids coatings unsuccessfully. The major problem encountered was very slow drying. This was a result of an improper level of catalysts in the coating.\nIM coatings for bridges using principally an epoxy mastics and silicone rubbers, that are not particularly relevant to aircraft topcoats, were evaluated. Cyclic salt-fog/freeze provided a relatively rapid method to differentiate coating performance in a short time period. Specific VOC content was not stated, but all the evaluated coatings were at or below 340 g/l. A xe2x80x9clow-VOCxe2x80x9d acrylic aliphatic polyurethane topcoat exhibited the best gloss retention.\nThe basic objectives of this invention are to produce a polymeric binder or matrix for a coating with extremely low, or zero VOC content that can be used as a top coat for many applications and in aircraft applications to achieve specific coating characteristics to prevent icing of aircraft wings. Icing on critical aircraft surfaces may create a condition which might impair the stability of the aircraft. The specific areas are referred to as xe2x80x9ccold-soakxe2x80x9d areas and some other areas on the xe2x80x9cleading edgesxe2x80x9d of the wings and engine nacelles. The present invention eliminates the adhesion of ice on these surfaces. Environmental icing due to weather is a related problem, but is not the direct problem concerning the present invention.\nIce will not adhere to the surface of certain polymer coatings with low surface energy such as Teflon. This is a consequence of the high contact angle between the water droplet and the surface that establishes a non-wetting surface. One objective of the present invention is to implement such coatings and a deposition process. Effective implementation will also result in a coating formulation and deposition process, with a very low VOC emission.\nCoatings formulations for prevention of icing problems includes the following properties:\nLow surface energies to prevent icing.\nAdhesive to aircraft surfaces.\nProtection of substrates from corrosion.\nResistance to jet fuel and hydraulic fluids.\nOther properties for coatings specifications.\nCoating materials are selected for low surface energy properties and general coating properties. The coating materials are polymerized fluoropolymers that possess good low temperature properties, e.g., do not embrittle at xe2x88x9245xc2x0 C., and do not soften at elevated temperatures of 90xc2x0 C.\nTwo parallel benefits of this approach may be achieved. The coating process is potentially adaptable for coating an entire aircraft or other commercial item. The combination of the coating process and the coating formulation reduces volatile organic compounds (VOC) well below the current Environmental Protection Agency (and California) limits for the forseeable future.\nFurther scope of applicability of the present invention will become apparent from the detailed description given hereinafter. However, it should be understood that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description."}
{"text": "The use of ink jet printers for printing information on a recording media is well established. Printers employed for this purpose may be grouped into those that continuously emit a stream of fluid droplets, and those that emit droplets only when corresponding information is to be printed. The former group is generally known as continuous inkjet printers and the latter as drop-on-demand inkjet printers. The general principles of operation of both of these groups of printers are very well recorded. Drop-on-demand inkjet printers have become the predominant type of printer for use in home computing systems, whereas continuous inkjet systems find major application in industrial and professional environments. Typically, continuous inkjet systems produce higher quality images at higher speeds than drop-on-demand systems.\nContinuous inkjet systems typically have a print head that incorporates a fluid supply system for fluid and a nozzle plate with one or more nozzles fed by the fluid supply. The fluid is jetted through the nozzle plate to form one or more thread-like streams of fluid from which corresponding streams of droplets are formed. Within each of the streams of droplets, some droplets are selected to be printed on a recording surface, while other droplets are selected not to be printed, and are consequently guttered. A gutter assembly is typically positioned downstream from the nozzle plate in the flight path of the droplets to be guttered.\nIn order to create the stream of droplets, a droplet generator is associated with the print head. The droplet generator stimulates the stream of fluid within and just beyond the print head, by a variety of mechanisms known in the art, at a frequency that forces continuous streams of fluid to be broken up into a series of droplets at a specific break-off point within the vicinity of the nozzle plate. In the simplest case, this stimulation is carried out at a fixed frequency that is calculated to be optimal for the particular fluid, and which matches a characteristic drop spacing of the fluid jet ejected from the nozzle orifice. The distance between successively formed droplets, S, is related to the jet velocity, v, and the stimulation frequency, f, by the relationship: v=fS. U.S. Pat. No. 3,596,275, issued to Sweet, discloses three types of fixed frequency generation of droplets with a constant velocity and mass for a continuous inkjet recorder. The first technique involves vibrating the nozzle itself. The second technique imposes a pressure variation on the fluid in the nozzle by means of a piezoelectric transducer placed typically within the cavity feeding the nozzle. A third technique involves exciting a fluid jet electrohydrodynamically (EHD) with an EHD droplet stimulation electrode.\nAdditionally, continuous inkjet systems employed in high quality printing operations typically require small closely spaced nozzles with highly uniform manufacturing tolerances. Fluid forced under pressure through these nozzles typically causes the ejection of small droplets, on the order of a few pico-liters in size, traveling at speeds from 10 to 50 meters per second. These droplets are generated at a rate ranging from tens to many hundreds of kilohertz. Small, closely spaced nozzles, with highly consistent geometry and placement can be constructed using micro-machining technologies such as those found in the semiconductor industry. Typically, nozzle channel plates produced by these techniques are typically made from materials such as silicon and other materials commonly employed in micromachining manufacture (MEMS). Multi-layer combinations of materials can be employed with different functional properties including electrical conductivity. Micro-machining technologies may include etching. Therefore through-holes can be etched in the nozzle plate substrate to produce the nozzles. These etching techniques may include wet chemical, inert plasma or chemically reactive plasma etching processes. The micro-machining methods employed to produce the nozzle channel plates may also be used to produce other structures in the print head. These other structures may include ink feed channels and ink reservoirs. Thus, an array of nozzle channels may be formed by etching through the surface of a substrate into a large recess or reservoir which itself is formed by etching from the other side of the substrate.\nFIG. 1 schematically illustrates a prior art conventional electrohydrodynamic (EHD) stimulation means used to excite a jet of conductive fluid into a stream of droplets. Fluid supply 10 contains conductive fluid 12 under pressure which forces ink through nozzle channel 20 in the form of a conductive fluid jet 22. Conductive fluid 12 is grounded or otherwise connected through an electrical pathway. A prior art droplet stimulation electrode 15 is approximately concentric with an exit orifice 21 of nozzle channel 20 as shown in cross-section in FIG. 1A. Droplet stimulation electrode 15 typically includes a conductive electrode structure 13 produced from a variety of conductive materials, including a surface metallization layer, or from one or more layers of a semiconductor substrate doped to achieve certain conductivity levels. Prior art conductive electrode structure 13 is electrically connected to a stimulation signal driver 17 that produces a potential waveform of chosen voltage amplitude, period and functional relationship with respect to time in accordance to a stimulation signal 19. In FIG. 1, an example of a stimulation signal 19 comprises a uni-polar square wave with a 50% duty cycle. The resulting EHD stimulation is a function of the square of field strength created at the surface of the conductive fluid 12 near exit orifice 21. The resulting EHD stimulation induces charge in the conductive fluid jet 22 and creates pressure variations along the jet. Conductive electrode structure 13 is covered by one or more insulating layers 24 which are necessary to isolate droplet stimulation electrode 15 from conductive fluid 12 in order to prevent field collapse, excessive current draw and/or resistive heating of conductive fluid 12. The conductive fluid 12 must be sufficiently conductive to allow charge to move through the fluid from the grounded fluid supply 10 in order to electrohydrodynamically stimulate conductive fluid jet 22 to form droplets that subsequently form at break-off point 26. Since conductive fluids are employed, a non-uniform distribution of charge cannot be supported in the fluid jet column outside of the stimulating electric field. The electrohydrodynamic stimulation effect occurs due to the momentary induction of charge in conductive fluid 12 at nozzle orifice 20 that creates the pressure variation in fluid jet 22. For a correctly chosen frequency of the stimulation signal 19, the perturbation arising from the pressure variations will grow on the conductive fluid jet 22 until break-off occurs at the break-off point 26.\nVarious means for distinguishing or characterizing printing droplets from non-printing droplets in the continuous stream of droplets have been described in the art. One commonly used practice is that of electrostatic charging and electrostatic deflecting of selected droplets as described in U.S. Pat. No. 1,941,001, issued to Hansell, and U.S. Pat. No. 3,373,437, issued to Sweet et al. In these patents, a charge electrode is positioned adjacent to the break-off point of fluid jet. Charge voltages are applied to this electrode thus generating an electric field in the region where droplets separate from the fluid. The function of the charge electrode is to selectively charge the droplets as they break off from the fluid jet.\nReferring back to FIG. 1, a typical prior art electrostatic droplet characterizing means includes charging electrode 30. Conductive fluid 12 is employed such that a current return path exists through the fluid supply 10 (e.g. through grounding). A charge is induced in a specific droplet under the influence of the field generated by charge electrode 30. This droplet charge is locked in on the droplet when it separates from the fluid jet 22. Charging electrode 30 is electrically connected to charge electrode driver 32. The charging electrode 30 is driven by a time varying voltage. The voltage attracts charge through conductive fluid 12 to the end of the fluid stream where it becomes locked-in or captured on charged droplets 34 once they break-off from the jet 22.\nA high level of conductivity of fluid 12 is required to effectively charge droplets formed in these prior art systems. Prior art inkjet print heads that employ electrostatic droplet characterizing means typically use conductive fluid 12 conductivities on the order of 5 mS/cm. These conductivity levels permit induction of sufficient charge on charged droplets 34 to allow downstream electrostatic deflection. The conductivity required for droplet charging is typically much greater than that for droplet stimulation. Typically, a conductive fluid suitable for charging can also be stimulated using EHD principles. The selective charging of the droplets in conventional electrostatic prior art inkjet systems allows each droplet to be characterized. That is, the conductive inks permit charges of varying levels and polarities to be selectively induced on the droplets such that they can be characterized for different purposes. Such purposes may include selectively characterizing each of the droplets to be used for printing or to not be used for printing.\nAgain referring to the prior art system shown in FIG. 1, a potential waveform produced by the charging electrode driver 32 will determine how the formed droplets will be characterized. The potential waveform will determine which of the formed droplets will be selected for printing and which of the formed droplets will not be selected for printing. Droplets in this example are characterized by charging as shown by charged droplets 34 and uncharged droplets 36. Since a specific droplet characterization is dependant upon whether that droplet is printed with or not, the potential waveform will typically be based at least in part on a print-data stream provided by one or more systems controllers (not shown). The print-data stream typically comprises instructions as to which of the specific droplets within the stream of droplets are to be printed with, or not printed with. The potential waveform will therefore vary in accordance with the image content of the specific image to be reproduced.\nAdditionally, the potential waveform may also be based on methods or schemes employed to improve various printing quality aspects such as the placement accuracy of droplets selected for printing. Guard drop schemes are an example of these methods. Guard drop schemes typically define a regular repeating pattern of specific droplets within the continuous stream of droplets. These specific droplets, which may be selected to print with if required by the print-data stream, are referred to as “print-selectable” droplets. The pattern is additionally arranged such that additional droplets separate the print-selectable droplets. These additional droplets cannot be printed with regardless of the print-data stream and are referred to as “non-print selectable” droplets. This is done so as to minimize unwanted electrostatic field effects between the successive print-selectable droplets. Guard drop schemes may be programmed into one or more systems controllers (not shown) and will therefore alter the potential waveform so as to define the print-selectable droplets. The voltage waveform will therefore characterize printing droplets from non-printing droplets by selectively charging individual droplets within the stream of droplets in accordance with the print data stream and any guard drop scheme that is employed.\nAgain referring to the prior art system shown in FIG. 1, electrostatic deflection plates 38 placed near the trajectory of the characterized droplets interact with charged droplets 34 by steering them according to their charge and the electric field between the plates. In this example, charged droplets 34 that are deflected by deflection plates 38 are collected on a gutter 40 while uncharged droplets 36 pass through substantially un-deflected and are deposited on a receiver surface 42. In other systems, this situation may be reversed with the deflected charged droplets being deposited on the receiver surface 42. In either case, further complications arise from the fact that the charging electrode driver 32 must be synchronized with stimulation signal driver 17 to ensure that optimum charge levels are transferred to droplets, thus ensuring accurate droplet printing or guttering as the architecture of the recorder may dictate. These synchronization constraints arise as result of charging or characterizing those conductive fluid droplets at a place and time separate from their stimulation. Although prior art electrostatic characterization and deflection systems are advantageous in that they permit large droplet deflection, they have the disadvantage that they have been used primarily only with conductive fluids, thus limiting the applications of these systems.\nA wide range of fluid properties is desirable in commercial inkjet applications. Jetted inks may be made with pigments or dyes suspended or dissolved in fluid mediums comprised of oils, solvents, polymers or water. These fluids typically have a large range of physical properties including viscosity, surface tension and conductivity. Some of these fluids are considered to be non-conductive fluids, and thus have insufficient levels of conductivity so as to be employed in continuous inkjet systems that rely on the selective electrostatic charging and deflection of conductive fluid droplets.\nVarious systems and methods for stimulating a non-conductive fluid medium to form a series of droplets and for characterizing the series of droplets to form “printing” droplets and “non-printing” droplets have been proposed. For example, U.S. Pat. No. 3,949,410, issued to Bassous et al., teaches use of a monolithic structure useful for the EHD stimulation of conductive fluid droplets in a jet stream emitted from a nozzle.\nU.S. Pat. No. 6,312,110, issued to Darty, and U.S. Pat. No. 6,154,226, issued to York et al., teach the construction of various inkjet print heads wherein droplets are not stimulated from a stream of non-conductive fluid. Rather, the print heads comprises EHD pumps within the print head nozzles themselves. Droplets are ejected from the fluid supply in a similar fashion to drop-on-demand printers.\nU.S. Pat. No. 4,190,844, issued to Taylor, teaches a use of a first pneumatic deflector for deflecting non-printing ink droplets towards a droplet catcher. A second pneumatic deflector either creates an “on-off” basis for line-at-a-time printing, or a continuous basis for character-by-character printing.\nU.S. Pat. No. 6,079,821, issued to Chwalek et al., teaches a use of asymmetric heaters to both create and deflect individual droplets formed in a continuous inkjet recorder. Deflection of the droplets occurs by the asymmetrical heating of the jetted stream.\nU.S. Pat. No. 4,123,760, issued to Hou, teaches the use of deflection electrodes upstream of a break-off point from which droplets are formed from a corresponding jetted fluid stream. Droplets produced by the stream are steered to different laterally separated printing locations by applying a cyclic differential charging signal to the deflection electrodes. This causes a deflection of the unbroken fluid stream which directs the droplets towards their desired printing positions.\nIt can be seen that there is a need to provide an apparatus and method of stimulating or forming a non-conductive fluid droplet or droplets from a jet of non-conductive fluid."}
{"text": "1. Field of the Invention\nThe present invention relates to a fuel cell system.\n2. Description of Related Art\nIn recent years, fuel cells, which generate electricity by using hydrogen (fuel gas) supplied to anodes and air containing oxygen (oxidizer gas) supplied to cathodes, has been developed and are expected as electric power sources for vehicles (e.g., fuel cell electric vehicles), etc.\nAs the electricity generation by such the fuel cell progresses, water vapor (water) is generated at the cathode of each MEA (Membrane Electrode Assembly) of the fuel cell, and part of the generated water permeates the MEA and then cross-leaks to the anode channel (fuel gas channel). Also, part of the cross-leaked water remains in the anode channel and adheres to the surface of the MEA, etc. The water adhering to the surface of the MEA obstructs the supply of hydrogen (fuel gas) to the anode (catalysts forming the anode). Consequently, the voltage (output) of the fuel cell can fall below a target voltage and electricity generation status of the fuel cell can degenerate to a poor condition.\nIn a technique proposed to solve the above problem, when the electricity generation status of a fuel cell has degenerated to a poor condition as above, gas and moisture remaining in the anode channel are purged (discharged) to the outside by temporarily raising the gas pressure in the anode channel and then briefly opening a purge valve (discharge valve) placed downstream of the anode channel (see JP 2008-112585 A)."}
{"text": "Spare tires on motor vehicles are rarely thought of by the vehicle user until such time that they are needed. When that happens, the vehicle user often finds out that the spare tire is flat or at least severely under inflated. Consequently, in order to be sure that a spare tire is fully inflated, the vehicle user must check the pressure therein from time to time, and this quite often requires removal of the tire from the trunk or moving it around in the trunk, either of which requires considerable effort. This problem is only exacerbated if the tire is secured to the trunk or if the trunk contains other things which must also be moved or removed during the process.\nU.S. Pat. No. 3,019,831 to Morrello attempts to solve the aforementioned problem by attaching a pressure hose to the valve stem on the spare tire inside of a trunk and has a pressure hose leading to a valve stem which extends to the outside of the vehicle body and is attached to the vehicle body so that the pressure can be checked at the same time the other tires on the vehicle are checked. It also permits inflation without opening the trunk of the car. While this approach to solving the problem is good in certain respects, it does require alteration to the body of the vehicle and also presents an unsightly valve stem sticking out of the side of the vehicle body. This valve stem extension can catch car wash brushes and can involve considerable expense in drilling through the body of a car and putting it into place.\nConsequently, there is a need for an apparatus for maintaining proper inflation of a spare tire without requiring modification to the body of a vehicle."}
{"text": "1. Field of the Invention\nThis invention relates to a novel method for the production of a halogen-containing phthalocyanine compound. More particularly, this invention relates to a method for the production of a halogen-containing phthalocyanine compound manifesting the absorption of heat ray and/or the absorption in a near infrared region, excelling in solubility in a solvent, and having excellent light-fastness.\nThe phthalocyanine compound which is obtained by the method contemplated by this invention excels in a heat ray-absorbing property and, therefore, is useful as a heat ray-absorbing dye. Since the phthalocyanine compound obtained by the method according to this invention absorbs light in a near infrared region having a wavelength in the range of 600-1000 nm, it can manifest excellent effects when it is used as a near infrared absorption dye and a near infrared sensitizer for writing or reading in an optical recording medium using a semiconductor laser, a liquid crystal display device, and an optical character reader, as a photothermal modifier for thermal transfer and a thermal paper, thermal stencil printing, as a near infrared absorption filter as for plasma display panel (PDP), as an asthenopia inhibitor, as a near infrared absorbing material for a photoconductive material, or as a color separation filter for an image pickup tube, as a color filter for liquid crystal display, as a selective absorption filter for a color braun tube, as a color toner, as a toner dye for flash fixing, as an ink jet ink, as an indelible bar code ink, also as a microorganism inactivating agent, as a photosensitive dye for oncotherapy, as a heat ray-shielding agent for automobiles and buildings, and as a discriminating agent for resin sorting.\n2. Description of Related Art\nThe needs for a near infrared absorbent dye have been mounting in consequence of the expansion of the field of the applications thereof. As regards the near infrared absorption dye to be used as a near infrared absorption dye and a near infrared sensitizer for writing or reading in an optical recording medium using a semiconductor laser, a liquid crystal display device, and an optical character reader, as a photothermal modifier for thermal transfer and a thermal paper, thermal stencil printing, as a near infrared absorption filter as for plasma display panel (PDP), as an asthenopia inhibitor, as a near infrared absorbing material for a photoconductive material, or as a color separation filter for an image pickup tube, as a color filter for liquid crystal display, as a selective absorption filter for a color braun tube, as a color toner, as a toner dye for flash fixing, as an ink jet ink, as an indelible bar code ink, also as a microorganism inactivating agent, as a photosensitive dye for oncotherapy, as a heat ray-shielding agent for automobiles and buildings, and as a discriminating agent for resin sorting, methods for the production of materials which satisfy wholly such characteristics as light-fastness, heat-resistance, and solubility (or compatibility with a resin) have been studied hitherto. None of the methods developed to date, however, has proved advantageous for commercial applications.\nThe methods for producing such phthalocyanine compounds have been disclosed in JP-A-2000-63693, JP-A-2000-169743, JP-B HEI-07(1995)-103318, and U.S. 2003/0234995 A1, for example.\nJP-A-2000-69693 discloses a method which comprises inducing the reaction of tetrafluorophthalonitrile with vanadium trichloride by using α-methyl naphthalene or benzonitrile solely while introducing an oxygen-containing gas such as a gas obtained by diluting air into the resultant reaction solution. Since the solvent, α-methyl naphthalene, to be used in this method is very expensive and since the reaction of this method necessitate the introduction of a hot oxygen-containing gas and, therefore, entails a possibility of the introduced gas entraining the expensive solvent and dispersing it outside the reaction system and threatening the danger of explosion, this method hardly deserves to be rated as favorable for commercialization. JP-A-2000-169743 discloses the reaction of tetrafluorophthalonitrile and vanadium trioxide in benzonitrile as a solvent in the presence of paratoluene sulfonic acid and calcium carbonate at 150° C. This official gazette recommends to use as an organic solvent in the reaction an inert substance having no reactivity with the starting materials and the reaction vessel and, for the purpose of preventing the reaction vessel from corrosion, to perform the reaction of a metal oxide and paratoluene sulfonic acid at relatively low temperature. The present inventors' review of this method, however, has revealed that this method is not appropriate because the phthalocyanine compound produced thereby as aimed at occasionally fails to manifest the performance as claimed.\nJP-B-HEI 07(1995)-103318 discloses a novel phthalocyanine compound and a method for the production thereof and claims to produce a fluorine-containing phthalocyanine compound by reacting a corresponding phthalonitrile compound with vanadium chloride in an inert solvent or an aprotic solvent. Though only ethylene glycol cited as a reaction solvent is an example of an inert solvent, no other alcohols cannot be found anywhere. The working examples cited therein include the synthesis of phthalocyanine in a current of nitrogen. This synthesis is found to use benzonitrile exclusively as an inert solvent. U.S. 2003/0234995 A1 discloses the reaction of a substituted phthalonitrile compound with vanadium chloride using benzonitrile and octanol under reflux. Since the reaction is effected without using an inert gas such as nitrogen, the yield is as low as 63.2 mol %."}
{"text": "(1) Field of the Invention\nThis invention relates to a thin film transistor (TFT) substrate and a method for fabricating such a TFT substrate and, more particularly, to a stagger type TFT substrate fabricated through a plurality of exposure processes using a plurality of kinds of masks and a method for fabricating such a stagger type TFT substrate.\n(2) Description of the Related Art\nIn recent years improvement in the performance of liquid crystal display units has been required greatly to obtain high brilliance, high intensity and high definition. In particular, the study and development of TFT substrates included in liquid crystal panels used in them are proceeding. In addition, today liquid crystal display units must be manufactured by efficient, more simplified processes from the viewpoint of corporate profits. Furthermore, consideration must be given to the environment on the earth in utilizing electric power and raw materials in these manufacturing processes.\nCurrently, TFT substrates are manufactured mainly by a reverse stagger system. TFT substrates are formed through a plurality of exposure processes by the use of a plurality of masks. In this case, at least five masks are used for forming a gate bus-line layer, an operation island layer, a drain bus-line layer, a protection film layer, and a pixel electrode layer on TFT substrates. That is to say, to manufacture reverse-stagger type TFT substrates, as many as five exposure processes must be performed. Therefore, to improve efficiency in the manufacture of such reverse-stagger type TFT substrates, a method, for example, for reducing the number of the exposure processes from five to three by the use of a half tone mask is proposed (see, for example, Japanese Unexamined Patent Publication No. 2001-311965).\nBy the way, there is another system, which is called a stagger system, for manufacturing TFT substrates. By adopting this system, TFT substrates which differ from reverse-stagger type TFT substrates in structure can be formed. Compared with reverse-stagger type TFT substrates, stagger type TFT substrates have, for example, the following advantages. Layers included in stagger type TFT substrates can be formed continuously. In addition, an interface treatment process is not necessary for manufacturing stagger type TFT substrates.\nConventionally, however, a plurality of exposure processes must be performed with five to eight masks to manufacture stagger type TFT substrates. Accordingly, manufacturing efficiency is low and it takes many a day to complete TFT substrates. Moreover, many (kinds of) masks are used and a plurality of exposure processes are performed. This increases the probability that a display defect will occur due to, for example, dust. As stated above, problems, such as manufacturing efficiency, a yield, and manufacturing costs, which will arise at the time of manufacturing stagger type TFT substrates remain unresolved."}
{"text": "The invention relates to a spring device for the motion drive of a movable component from a rest position into a displaced position. The inventive spring device includes a preloaded spring and a gas spring. The movable component is loaded into the displaced position by the preloaded spring. A pressurized gas-filled cylinder of the gas spring include a first end closed by a first end plate and a second end closed by a second end plate and the interior of which is subdivided by an axially displaceable piston into a first working chamber and a second working chamber. The first working chamber and the second working chamber are connected to each other via a restrictor, and the piston has a piston rod which is led through the first working chamber and is led to the outside through the first end plate in a sealed manner, the piston rod having a free end being fixed to the movable component or a stationary component.\nIn such spring devices, the extension force is composed of a sum of the force of the preloaded spring and the force component of the gas spring that is present in the extension direction, which is primarily intended to effect damping of the extension movement of the piston rod. The damping force depends on the pressure of the gas in the cylinder. As the pressure rises, the damping force also rises.\nAs resistances of the components to be moved increase, the forces of the spring must also increase. This alone requires a high expenditure of force for the manual movement of the movable component out of the displaced position into the rest position. If increased damping of the movement from the rest position into the displaced position is also intended to be provided, the increased damping further increases the required expenditure of force."}
{"text": "Within the past few years special types of light weight fracture braces have been used for the treatment of tibia fractures after the initial period of treatment in a cast of plaster of paris. The cast is used initially to immobilize the fracture. The plaster cast is uncomfortable because it is heavy and is not removable during its long period of use. It also limits the mobility of the patient. After the plaster cast is removed the special light weight fracture braces can be worn to increase patient mobility while providing the support necessary to prevent twisting, or other undue stress on the tibia during the healing process. These fracture braces are light in weight because they are commonly made of plastic; and they are usually designed to be removed, adjusted, and reused on the patient. Some of these light weight fracture braces include a foot plate attached to the bottom of a tibia brace in the form of a jacket or the like that wraps around the tibia. The foot plate is secured to the tibia supporting jacket through flexion joints that permit limited rotational motion of the ankle. This brace allow the patient a limited amount of mobility while the brace is worn. Thus, the light weight fracture braces provide a substantial benefit to the patient during the healing process when compared with a plaster cast.\nProblems have been experienced with prior light weight fracture braces. The supporting jackets of these braces are commonly applied from rear to front of the patient's leg. This allows the patient's foot to be held in a normal position while the jacket is being placed on the patient's leg. On the other hand, if the jacket is applied from front to rear, the prior braces require too much movement of the patient's foot. It is best if the patient's foot is not moved from its normal position when the brace is applied. The braces applied from rear to front use the fatty tissue of the calf in the back of the tibia for attaching the brace to the leg. The jacket has an opening at the front so that the main portion of the jacket can be wrapped around the calf, often while the front opening is adjusted at the front of the tibia. It is difficult to make a brace that conforms well to the size and shape of the fatty tissue of the calf for all patients while ensuring that the size of the jacket is easily adjusted at the front of the tibia. In many cases the point of adjustment in front of the tibia applies undue pressure to the bony prominence along the front of the patient's tibia. Tibia supporting jackets made of hard plastic also are especially uncomfortable when tightened against the bony front portion of the patient's tibia. Another problem is that many of these prior art fracture braces are not always comfortable when the tibia supporting jacket is tightened. They are not designed to ensure that torque is uniformly applied to the patient's leg when the straps on the tibia supporting jacket are wrapped around the jacket and tightened. The tibia supporting jacket is often crumpled or too bulky when tightened. In either case the brace can be highly uncomfortable for the patent.\nThis invention provides an orthosis that overcomes the problems associated with the prior art light weight fracture braces. The tibia supporting jacket of the orthosis can be applied from front to rear of the patient's tibia. The adjustment is in the rear of the jacket, along the fatty tissue of the patient's calf. The front of the jacket thus can provide continuous padding along the bony prominence along the front of the patient's tibia, while the adjustment at the rear of the jacket along the patient's calf is much more comfortable than the prior art braces in which the jacket is applied from rear to front. The orthosis also can be applied from front to rear while leaving the patient's foot in the normal position. In addition, the tibia supporting jacket can be applied and tightened without torquing the patient's leg one way or the other. The jacket is held in a neutral position around the tibia when tightened to its fullest extent, and the jacket also folds together and is tightened in such a way the the material does not crumple or become bulky. The result is an orthosis which is comfortable to wear while still providing the required amount of support. In addition to providing the support necessary for fractures of the tibia, the orthosis also can be used as a means of support for ankle sprains."}
{"text": "The present invention relates to fuel gas burners for cooking appliances and particularly burners of the type employed in cooktop or rangetop applications where a receptacle or cooking vessel is seated on the surface of the burner for heating of the foodstuffs or liquid within the vessel. Cooktop burners are typically ignited by the user opening a rotary valve in the supply line to provide a flow of the fuel gas to the burner whereupon a set of switch contacts are simultaneously closed for electrically energizing an igniter having an electrode disposed to provide a spark in the stream of fuel air mixture emanating from a port in the burner. If an alternating current voltage is employed for the spark ignitor, upon ignition of the fuel air mixture and the presence of flame about the ignitor electrode, the phenomenon of flame rectification occurs; and, the change in the current may be electrically detected as an indication or proof of the presence of flame. This technique has been widely employed for combining the function of the ignitor with that of a flame sensor and providing electrical circuitry which could respond to the change in alternating current to turn off the sparking voltage to the ignitor. It is also known to provide circuitry which, upon the loss of flame, electrically detects the change of a current in the electrode and reenergizes the ignitor spark voltage automatically. However, if transient air currents extinguish the flame about an annular plural port burner on only a portion of the periphery, the flame sensor may not be able to determine whether the flame has been totally extinguished and an annoying reenergization of the ignitor occurs. The condition may also occur where variations in the line pressure of the fuel gas cause major fluctuations in the flame.\nThus, it has long been desired to provide a way or means of preventing flame loss in the region of the flame sensing ignitor when flame is being sustained in other regions of the burner and to generally stabilize the flow from the flame generating ports in the burner. It has further been desired to improve the effectiveness of a spark ignitor for a cooktop burner and to provide such functions in a burner which is sufficiently low in manufacturing cost to remain competitive in the high-volume domestic appliance marketplace."}
{"text": "The invention relates to a method of improving a surface of a semiconductor substrate, wherein the surface at least partially includes silicon.\nIn semiconductor device production, it is more and more important to provide semiconductor substrates of a very high quality. Defects of semiconductor substrates can be of very different origin and may occur in the bulk material of wafers or layers or on the surface of a structure. Deficient wafers, such as wafers or layers with holes or scratches on their surface or with oxide precipitates or so-called “HF-defects”, which are present in or on a wafer and will be apparent by an HF-etch step, are mostly not suitable for further use.\nTo improve the surface characteristic of a defective wafer, a wafer treatment of a wafer surface such as an etching step or a chemical mechanical polishing (“CMP”) step can be used to remove or to reduce the number or the size of defects at or near the wafer surface. Typical etchants are halogen bearing compounds such as HCl, HBr, HI, HF, and others. The etchant can also be a fluorine bearing compound such as SF6, or CxFx. Moreover or in addition, it is possible to treat a defect containing wafer thermally, preferably in a hydrogen bearing environment, to smooth it and to diminish its defects. The thermal treatment can be performed in a furnace or in a tool for rapid thermal processing (“RTP”). According to another approach disclosed in U.S. Pat. No. 6,287,941 B1, a defective wafer such as a cleaved film can be subjected to a combination of etching and deposition at very high temperature using a combination of etchant and deposition gases to result in a better surface quality.\nAlthough such methods lead, in the first instance, to a superficial improvement of the surface condition of a defective wafer by smoothing, abrasion or defect covering of the respective wafer, the known methods are mostly very laborious and the corresponding defects cannot really be repaired. Thus, there remains a need to process defective wafers to increase surface quality."}
{"text": "Known commercial mobile communication systems typically include a plurality of fixed base stations arranged in patterns whereby each base station transmits and receives over a plurality of frequencies. A mobile station within range of the base station can communicate with the external world (e.g., via the Public Switched Telephone Network (“PSTN”)) through the base station using the frequencies. The area surrounding a base station in which mobile stations communicate with that base station is often referred to as a cell, with the base station generally positioned toward the center of the cell. Examples of known commercial mobile communications systems having cells include cellular communications systems, Personal Communications Systems (“PCS”), Global System for Mobile communication (“GSM”) systems, IS-136/Digital-American Mobile Phone systems (hereinafter “IS-136” or “D-AMPS”), and so on.\nIn an IS-136 system, a mobile station can communicate with the base station via a carrier frequency pair that includes two different (but paired) frequencies. The first frequency of the pair is the downlink (or forward) frequency where information is transmitted from the base station to the mobile station, and the second frequency of the pair is the uplink (or reverse) frequency where information is transmitted from the mobile station to the base station. Each carrier frequency pair is often referred to as a carrier or a channel, although the term channel is also used in different ways when a carrier can carry multiple channels (e.g., time-division multiple access (“TDMA”) channels, code-division multiple access (“CDMA”) channels, and so on). An IS-136 system can have 416 carriers, of which 395 carriers are available to carry voice traffic between a mobile station and a base station.\nThe carriers used by a base station are separated from one another in frequency to minimize interference. A cell's carriers are carefully selected so that adjoining cells do not transmit or receive on the same carrier frequencies. A mobile system operator can allocate to a base station a set of carriers with frequencies that are each separated from the next carrier by an integral number. For example, FIG. 1 shows a known frequency reuse pattern for base station cells, where each cell is assigned a set of carriers. Each cell can be allocated one of seven sets of carriers, where a cluster of seven cells as a whole is allocated all of the carriers. Thus, the frequency reuse pattern illustrated in FIG. 1 is typically referred to as having n=7 clusters. The cells are arranged and frequency sets can be allocated by assigning a first carrier set (e.g., carrier set 1) to a central cell 111 of a first cluster, and then assigning different carrier sets (e.g., carrier sets 2–7) to the cells of the first cluster surrounding that central cell. Thus, cell 111 can have carrier set 1, cell 112 can have carrier set 2, cell 113 can have carrier set 3, cell 114 can have carrier set 4, cell 115 can have carrier set 5, cell 116 can have carrier set 6, and cell 117 can have carrier set 7. Each of the carrier sets are also respectively allocated to the cells 141–147 of a fourth cluster adjacent to cells 111–117 of the first cluster. Portions of other adjacent clusters—such as cells 125 and 126 of a second cluster, cells 131 and 134–137 of a third cluster, and so on—are also illustrated.\nThe assignment of carriers to carrier sets, and the assignment of carrier sets to cells, can be based on the number of different carrier sets (e.g., seven, four, and three carrier sets) and the number of available carriers. An IS-136 system having 395 voice carriers and using a frequency reuse pattern illustrated in FIG. 1 can have approximately 57 carriers per carrier set and cell. With seven different carrier sets, carrier set 1 can include carriers 1, 8, 15, 22, 29, 36 and so on; carrier set 2 can include carriers 2, 9, 16, 23, 30, 37 and so on; carrier set 3 can include carriers 3, 10, 17, 24, 31, 38 on so on; and so forth with respect to carrier sets 4–7.\nFIGS. 2 and 3 illustrate other known frequency reuse patterns. In particular, FIG. 2 illustrates a frequency reuse plan having n=3 clusters. A first cluster can have three cells 211–213. Six clusters—such as a second cluster having cells 221–223, a third cluster having cells 231–233, a fourth cluster having cells 241–243, and so on—can be located adjacent the first cluster. Each cell of each cluster can be allocated a third of the available system carriers. FIG. 3 illustrates a frequency reuse plan having n=4 clusters. A first cluster can have four cells 411–114. Six clusters—such as a second cluster having cells 321–324, a third cluster having cells 331–334, a fourth cluster having cells 341–344, and so on—can be located adjacent the first cluster. Each cell of each cluster can be allocated a fourth of the available system carriers.\nTo allow a mobile station to transmit and receive communications as the mobile station moves from one cell to another, each cell is normally positioned with its area of coverage overlapping the areas of coverage of a number of adjacent and surrounding cells. As a mobile station moves from an area covered by a first base station to an area covered by another base station, mobile station communications (e.g., a voice call, a data link, etc) are transferred from the first base station to the other base station in an area where the coverage from the two cells overlaps. The transfer of a mobile station from communicating with one base station to communicating with a second base station is typically called hand-off.\nA cell can have at least two types of radio coverage. A first type of cell radio coverage is omnidirectional (i.e., azimuthally), where the cell has an antenna set that can communicate with mobile stations via each carrier of the carrier set allocated to the cell. A second type of cell radio coverage is sectored. FIG. 4 shows an illustration of a sectored cell. Cell 411 includes a plurality of sectors, including sectors 401, 402 and 403. Sectors are often referred to as an alpha sector, a beta sector, and a gamma sector. Cells are typically divided into three sectors, with each sector having an antenna set that covers a 120° sector. In a cell having three sectors, each sector antenna set can communicate with mobile stations via one-third of the carriers of the carrier set allocated to the cell so that each sector communicates over different carriers as compared to the other sectors of the cell.\nNotwithstanding the use of frequency reuse patterns, interference between like carriers of different cells can occur. For example, referring again to FIG. 1, even though cell 131 is a knight's move away from cell 111 (i.e., cell 131 is up two cells and over one cell from cell 111), there can be interference between the carriers of cell 111 and cell 131. For example, within portions of cells 113, 137, and 136, there can be interference between a carrier 1 of cell 111 and a carrier 1 of cell 131. Such interference is typically called co-channel interference.\nCo-channel interference can be caused by antenna patterns, power levels, carrier scattering, and wave diffraction that differ from cell to cell. Buildings, structures, mountains, foliage, and other physical objects can cause carrier signal strength to vary over the area covered by a cell. As a result, the boundaries (i.e., edges) at which the signal strength of a carrier falls below a level sufficient to support communications with a mobile station can vary widely from cell to cell. Thus, cells adjacent one another do not typically form anything like the precise geometric patterns illustrated in FIGS. 1–3. Cell coverages, however, must overlap to allow mobile stations to be handed-off between cells, and such overlapping, among other factors, can lead to co-channel interference.\nIn an IS-136 system, mobile stations are instructed to measure the signal strengths of various carriers and report the measured signal strengths to the mobile system. For example, referring again to FIG. 1, as a mobile station in communication with the base station of cell 111 moves through cell 111 toward cells 112 and 113, the mobile station can be instructed to measure the signal strengths of certain carriers of cells 111, 112, and 113 and report the measured carrier signal strengths to the mobile system via the base station of cell 111. When the signal strength reported by the mobile station with respect to the cell 111 carrier drops below a certain threshold (e.g., as the mobile station approaches the intersection of cells 111, 112, and 113), the mobile system will pick one carrier of the carriers measured and reported by the mobile station and instruct the mobile station to use that carrier for communications (e.g., instruct the mobile station to begin communicating with the base station of cell 113 or cell 112 via the appropriate carrier). In known IS-136 systems, mobile stations can monitor and report the carrier strengths of neighboring surrounding cells.\nMobile system operators generate frequency reuse plans to, among other things, reasonably minimize co-channel interference and reasonably maximize the likelihood that mobile stations will be successfully handed-off to a next cell as it moves away from its current cell. A first method of generating a mobile system frequency reuse plan is to use a frequency reuse pattern as illustrated in FIGS. 1–3. The efficiency of a frequency reuse plan can be increased by modifying the frequency reuse plan based on knowledge (subjective and/or objective) of the terrain covered by the frequency reuse plan. For example, FIG. 5 illustrates a frequency reuse pattern that has been modified based on terrain characteristics. A first cluster includes cells 511–517. A mountain range 501 abuts the edge of the first cluster at the exterior edges of cells 513 and 514. The mountain range 501 attenuates the carrier signals transmitted by cells 512, 511, and 515. Thus, cells 522, 521, and 525 can use the same carrier sets used by cells 512, 511, and 515 with a reasonable minimization of co-channel interference. By reusing the carrier sets in a more compact manner, the frequency reuse plan illustrated in FIG. 5 is more efficient than the frequency reuse pattern illustrated in FIG. 1. A more efficient frequency reuse plan allows for greater system utilization (e.g., more mobile stations can be supported).\nFrequency reuse plans can also be based on predictive methods using computer modeling. A computer model can predict carrier propagation areas based on antenna height, transmitter power, terrain characteristics, and so forth. Measured carrier data can also be used to create and modify frequency reuse plans. In an IS-136 system, mobile stations report received carrier strengths to the mobile switch coupled to the base stations. The reported carrier strength data can be used to determine how far carriers propagate.\nCarrier propagation and co-channel interference can also be measured by receivers that measure received carrier strength as they are driven throughout areas of the mobile system during a so-called “drive test.” For example, during a drive test a specific test carrier is transmitted at each cell or sector of a cell involved in the interference testing. A scanning receiver is driven over the roads, highways and traveled byways of the system. The scanning receiver scans and measures the strength of the test carrier signal transmitted by each cell at the points of possible interference, and location determination equipment (e.g., a Global Positioning System (“GPS”) unit, a Loran unit, etc.) records the position of the scanning receiving. These strength measurements are then plotted and the expected interference points from different cells may be viewed graphically to determine whether sufficient interference exists to change the channel sets assigned to a particular area. This method of performing a drive test is often referred to as a “key-up” drive test because the test carrier is continuously “keyed-up” at each cell so as to be measured. A test carrier does not carry subscriber communications.\nU.S. Pat. No. 5,926,762 (“the '762 patent”) describes another type of drive test in which a unique test carrier at each cell site is transmitted such that each cell site is transmitting a different test carrier. A scanning receiver is driven over the roads, highways and traveled byways of the system to measure the strength (typically the received signal power) of each test carrier transmitted by each of the cell sites while location determination equipment records the position of the scanning receiver. According to the '762 patent, transmitting a different test carrier at each cell eliminates interference that can complicate strength measurements when a single carrier is keyed-up at multiple cells for a drive test.\nThese known methods of performing drive tests to measure carrier strengths and predict co-channel interference require test carriers to be keyed-up to continuously transmit. Whether a single test carrier is keyed-up at a plurality of cells, or different test carriers are keyed-up at different cells, each method requires keying-up a test carrier. When a test carrier is keyed-up, it is not available to carry subscriber communications (e.g., voice traffic), and system capacity is diminished. Thus, key-up drive tests are typically conducted during the evening when demand for system capacity is lowest. In view of the foregoing, it can be appreciated that a substantial need exists for systems and methods that can advantageously provide for determining mobile communication system telephone carrier frequency propagation characteristics."}
{"text": "In moving picture coding processing, in general, the amount of information is reduced by utilizing redundancy in the spatial direction and the temporal direction which moving pictures have. Here, in general, transform to a frequency domain is used as a method utilizing redundancy in the spatial direction. Further, inter-picture prediction (hereinafter, referred to as “inter prediction”) coding processing is used as a method utilizing redundancy in the temporal direction. In inter prediction coding processing, when a picture is coded, a coded picture that appears before or after a current picture to be coded in the display time order is used as a reference picture. A motion vector is derived by performing motion detection on the current picture relative to the reference picture. Then, redundancy in the temporal direction is eliminated by calculating a difference between image data of the current picture and predicted image data obtained by motion compensation based on the derived motion vector."}
{"text": "1. Field of the Disclosure\nThe disclosure relates to a sense amplifier used in a semiconductor device, and particularly relates to a suitable sense amplifier in a semiconductor device that has a variable resistance memory cell, and to a data processing system.\n2. Description of Related Art\nConventional memory cells are known that store information based on the size of a resistance value or the “on” current of a transistor. This type of memory cell generally has relatively high resistance values ranging from 10 kΩ to several hundred kilohms (or kilo-ohm) even in a low memory state, and sense amplification is therefore usually performed using a highly sensitive differential current sense amplifier (see Japanese Patent Application Laid-Open No. 2004-39231)."}
{"text": "In general, a golf club includes a club face with grooves in order to increase the frictional force at impact between the club face and the ball. The United States Golf Association (USGA) publishes and maintains the Rules of Golf, which govern professional golf in the United States. Appendix II to the USGA Rules provides several limitations on the grooves for golf clubs, including limitations on symmetry, width, depth, edge radius, and relative distance between grooves. The Royal and Ancient Golf Club of St. Andrews, which is the governing authority for the rules of professional golf outside the United States, provides similar limitations to golf club design.\nMoreover, an accurate hitting stoke is accomplished through a variety of subjective, as well as objective, golf club features. For example, many people subjectively associate smaller grooves with decreased ball spin at impact and with shorter distances after impact.\nFor simplicity and clarity of illustration, the drawing figures illustrate the general manner of construction, and descriptions and details of well-known features and techniques may be omitted to avoid unnecessarily obscuring the invention. Additionally, elements in the drawing figures are not necessarily drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help improve understanding of embodiments of the present invention. The same reference numerals in different figures denote the same elements.\nThe terms “first,” “second,” “third,” “fourth,” and the like in the description and in the claims, if any, are used for distinguishing between similar elements and not necessarily for describing a particular sequential or chronological order. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments described herein are, for example, capable of operation in sequences other than those illustrated or otherwise described herein. Furthermore, the terms “include,” and “have,” and any variations thereof, are intended to cover a non-exclusive inclusion, such that a process, method, system, article, device, or apparatus that comprises a list of elements is not necessarily limited to those elements, but may include other elements not expressly listed or inherent to such process, method, system, article, device, or apparatus.\nThe terms “left,” “right,” “front,” “back,” “top,” “bottom,” “over,” “under,” and the like in the description and in the claims, if any, are used for descriptive purposes and not necessarily for describing permanent relative positions. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the invention described herein are, for example, capable of operation in other orientations than those illustrated or otherwise described herein.\nThe terms “couple,” “coupled,” “couples,” “coupling,” and the like should be broadly understood and refer to connecting two or more elements or signals, electrically, mechanically and/or otherwise."}
{"text": "Recently, non-volatile semiconductor flash memories have come into wide use as data storage memories with excellent portability. The per-bit price of those flash memories has been rapidly dropping just from their miniaturization. This per-bit price reduction has actually been achieved by the element structure improvement or employment of multi-bit storage systems with respect to those flash memories.\nTypical methods for forming memory arrays of large capacity flash memories used for files are NAND type and AND type. In the NAND type memory, memory cells are connected serially. In the AND type memory, memory cells are connected in parallel. The AND type memory in which memory cells are disposed in parallel is usually considered to be suitable for multi-bit storage operations, since it enables controlling of the number of electrons stored in a floating gate. Additionally, the AND type memory employs a hot electron writing method, so that its writing is fast. The NAND type is disclosed in “IEEE International Electron Devices Meeting” (pp. 775-778, 2000)” by F. Arai et al., while the AND type is disclosed in “IEEE International Electron Devices Meeting (pp. 29-32, 2001)” by T. Kobayashi et al.\nThe official gazette of JP-A 156275/2001 illustrates a non-volatile memory technique that achieves both requirements of an array configuration in which memory cells are connected in parallel and a small memory cell region. This gazette further illustrates how to use each inversion layer formed on a semiconductor substrate located under an assist gate as a line. Also illustrated by the official gazette of JP-A No.2001-326288 is a technique for configuring a memory cell array at narrow word line pitches to achieve high density disposition of memory cells.\nAs described above, the AND type flash memory, which employs the hot electron writing technique, is fast in writing. Because the hot electron writing method employs source side injection, the method is also considered to be suitable for simultaneous writing in many memory cells. Additionally, because memory cells in an array are connected in parallel, each memory cell is not affected by the information stored in other adjacent memory cells so easily. This is why the AND type flash memory is also considered to be suitable for multi-bit storage per cell.\nIn spite of such advantages, the AND type flash memory continues to present difficulties. Because the AND type flash memory has an array structure in which diffusion layers are disposed in parallel, it is difficult to reduce the line pitches that are parallel to data lines due to the spread of the diffusion layers or existence of isolation regions. To solve this problem, a method for using inversion layers formed under the electrodes disposed in parallel to the data lines as local data lines may enable the subject AND type flash memory to operate without diffusion layers to be formed by impurity injection. This method is illustrated in the official gazette of JP-A No. 156275/2001.\nHowever, each inversion layer usually has a resistance higher than that of the diffusion layer formed by means of high density impurity injection into the object semiconductor substrate. This is why the local data line resistance is different among places in the memory array, so that as the voltage falls, the potential to be applied to each target memory cell changes and the writing characteristic differs among memory cells significantly. This problem is accentuated as local data lines become longer. Another problem to arise from the employment of the above described memory structure is that if the flash memory is structured so that local data lines are connected to a global data line at a short distance through a switch simply, the number of memory cells per local data line is reduced and the area penalty of a selected transistor portion increases."}
{"text": "1. Field of the Invention\nThis invention relates to an applicator assembly including an applicator roll and a secondary or supplementary applicator element both mounted on a housing and extending outwardly from an open face thereof when the applicator roll is designed to paint a relatively large surface area such as walls or surfaces and the secondary applicator is designed to apply paint to the junctions of such wall or ceiling surfaces.\n2. Description of the Prior Art\nNumerous devices have been introduced into the prior art for the purpose of transferring or spreading paint or like material unto a given surface. In addition to the wide variety of designs of simple or basic paint brushes, paint applicator designs include relatively complex structures.\nPrior art paint applicators include an applicator roll which can rapidly and effectively apply a paint over relatively large surface areas in an efficient manner. Such rolling of paint unto a surface has been found to be much quicker than \"stroking\" with a conventional paint brush. Where such applicator rolls are allowed to be used such as on large, relatively flat planar surfaces the rolls have been found to be more efficient than a brush and much less costly especially when considering the man or labor hours involved.\nHowever, certain disadvantages have become associated with the use of applicator rolls. The overall painting with this type of applicator roll structure has frequently been found to be messy, often times resulting in the waste and loss of paint. In the prior art, applying paint to the roller surface usually includes a roll tray having an inclined base at least partially filled with paint in which the applicator roll is reciprocally moved. Paint is transferred on to the outer surface of the roll and then the roll is of course movably applied over the surface being painted. In addition to the loss of paint through waste, there is of course the requirement for relatively additional, specifically constructed containers in the form of a tray requiring the applicator roll to be constantly \"dipped\" into the tray. The applicator roll has received wide acceptance primarily due to its rapid even spreading of paint onto a given surface. However, the industry is desires of overcoming recognized disadvantages by providing a more efficient structure which accomplishes the overall result of painting in the same accepted manner but which is more efficient, less messy but yet is easy to operate and maintain. The following U.S. and foreign patents are representative of prior art attempts to overcome the problems set forth above and long recognized as existing in the prior art. Such patents include U.S. Pat. Nos. to Boyle; 375,919 Peterson, 356,695; Fernandez, 1,376,195; Sporer, 1,461,947; Rufo, 2,307,858; Tucker, 2,538,542; 2,548,,653; Leverock, 2,583,432; Ballard et al 3,403,960; Chrun, 2,746,071; Pedro, 2,928,1113; Cassidy, 3,193,868; Clark et al, 3,231,151: Schultz, 3,274,637; Ellis, 3,690,779; Ogenibene, 3,721,502; Bradshaw, 3,809,484; Hansen, 3,825,970; Linton, 3,925,927; Spransy, 4,012,151; Rearai, 4,059,358; Gamacher, 4,129,391; Garcia, 4,140,410; Miller, 4,222,678; British Patent Nos. 14,024, 887,294 and 480,837; Swiss Patent No. 421,769; and German Patent No. 2,447,848."}
{"text": "The present invention relates to electronic circuits, and more particularly to storage elements used in sequential logic circuits.\nFlip-flops are widely used in electronics circuits to store data. FIG. 1 is a block diagram of a flip-flop 10 commonly referred to as master-slave flip-flop. Master-slave flip-flop 10 includes a master D-latch 12, a slave D-latch 14, and an inverter 16. When clock signal CLK is at a high logic level (high), slave latch 14 is enabled, i.e., samples the data, and its output Q is the same as the master latch 12 output Y. When clock signal CLK is high, master latch 12 is disabled, i.e., latches (holds) its data. When signal CLK is at a low logic level (low), master latch 12 is enabled and the data in the external D input is transferred to the master latch 12 output Y. When signal CLK is low, slave latch 14 is disabled, holds its data, and thus any changes in the external D input changes the master latch output Y and cannot change the slave output Y. When signal CLK returns to high, master latch 12 is in the latched mode and is isolated from the D input. At the same time, slave latch 14 is enabled and the value of Y is transferred to the output Q.\nA D-latch may be formed using transmission gates and inverters, as shown in FIG. 2. Input CLK controls the two transmission gates 22, 24. When CLK is low, transmission gate 22 has a closed path and transmission gate 24 has an open path, therefore, D-latch 20 is in a sampling mode. This cause output Q to be the inverse of input D and output Q to be the same as input D. When CLK is high, transmission gate 24 has a closed path and transmission gate 22 has an open path, therefore, D-latch 20 is in a holding mode; this causes outputs Q and Q to hold their previous values.\nA need continues to exist for a flip-flop that can operate at higher data rates for any given clock frequency."}
{"text": "1. Field of the Invention\nThe present invention relates to a connector engagement detecting apparatus which has a means to determine whether or not a pair of mating connectors used for connection of automotive wiring harnesses are normally joined together.\n2. Prior Art\nReferring to FIGS. 6 and 7, one of mating connector housings a is formed with a contact accommodating chamber d in which a pair of electric contacts b, c are inserted in non-contacting condition. The other mating connector housing e has a drive piece f, formed as a resilient cantilever, whose free end f.sub.1 forces the lower contact c upward into contact with the upper contact b. The connector housing a also has an interfering projection g in front of the electric contact c, which, when the paired connector housings fail to be connected normally, abuts against the free end f.sub.1 of the drive piece f, deflecting it to block the electric contacts b, c from coming into forced contact with each other. When the mating connector housings are completely connected together, the interfering projection g is received into a recess f.sub.2 allowing the drive piece f to move from a position indicated by a broken line in FIG. 7b to a position of a solid line, which in turn causes the contact c to engage with the contact b to complete a detection circuit.\nIn the above-mentioned prior art, since the dedicated chamber d for accommodating the detecting electric contacts b, c is necessary, the connector housing becomes complex in shape, making the resin molding process correspondingly more difficult. Moreover, the drive piece f made of resin material may undergo thermal deformation from ambient heat generated during service. In that case, the driving force acting on the electric contact c decreases, degrading the reliability of electric conduction through the electric contacts b and c."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a server system, and in particular, to a server rack system.\n2. Description of Related Art\nIn recent years, a computer server gradually develops from a conventional single server to a rack server, where several servers are placed in a rack. Since a large number of servers in the rack server exist, management and control for the servers become the critical technology for the rack server. However, in the existing rack server, due to lack of a solution for efficiently managing and controlling the servers, the development of the rack server is seriously hindered."}
{"text": "1. Field of the Invention\nThis invention relates to a photomultiplier. This invention particularly relates to a long photomultiplier having a cylindrical main body with an elongated light receiving face that extends along the longitudinal direction of the main body.\n2. Description of the Prior Art\nWhen certain kinds of phosphors are exposed to radiation such as X-rays, .alpha.-rays, .beta.-rays, .gamma.-rays, cathode rays or ultraviolet rays, they store part of the energy of the radiation. Then, when the phosphor which has been exposed to the radiation is exposed to stimulating rays such as visible light, light is emitted by the phosphor in proportion to the amount of energy stored during exposure to the radiation. A phosphor exhibiting such properties is referred to as a stimulable phosphor.\nAs disclosed in U.S. Pat Nos. 4,258,264, 4,276,473, 4,315,318 and 4,387,428 and Japanese Unexamined Patent Publication No. 56(1981)-11395, it has been proposed to use stimulable phosphors in radiation image recording and reproducing systems. Specifically, a sheet provided with a layer of the stimulable phosphor (hereinafter referred to as a stimulable phosphor sheet) is first exposed to radiation which has passed through an object such as the human body in order to store a radiation image of the object thereon, and is then exposed to stimulating rays, such as a laser beam, which cause it to emit light in proportion to the amount of energy stored during exposure to the radiation. The light emitted by the stimulable phosphor sheet upon stimulation thereof is photoelectrically detected and converted to an electric image signal, which is processed as desired to reproduce a visible image having an improved image quality, which allows the visible image to be used in making efficient and accurate diagnoses of illnesses.\nIn general, radiation image read-out apparatuses used in the aforesaid radiation image recording and reproducing systems are constituted of a main scanning means for scanning a stimulable phosphor sheet, on which a radiation image has been stored, with stimulating rays in a main scanning direction, i.e. along main scanning line, a sub-scanning means for moving the stimulable phosphor sheet with respect to the stimulating rays in a sub-scanning direction approximately normal to the main scanning direction, and a photo detecting means for detecting light emitted by the stimulable phosphor sheet in proportion to the amount of energy stored during exposure to radiation.\nRecently, a novel photo detecting means which utilizes a long photomultiplier was proposed in, for example, Japanese Unexamined Patent Publication No. 62(1987)-16666. The disclosed long photo-multiplier is provided with a cylindrical main body having a light receiving face which extends along the main scanning line on the stimulable phosphor sheet. A photocathode is provided on an inner surface of the main body along the light receiving face. Also, in general, a light guide member is located so that it is in close contact with the light receiving face and so that it extends along the light receiving face and guides the light emitted by the stimulable phosphor sheet toward the photocathode. With the long photomultiplier, non-directional light emitted by the stimulable phosphor sheet is guided by the light guide member toward the photocathode. When exposed to light emitted by the stimulable phosphor sheet, the photocathode generates photoelectrons, which are sequentially multiplied by the secondary electron emission effects of dynodes.\nWith the aforesaid long photomultiplier, light emitted by every portion of the stimulable phosphor sheet in the main scanning direction can be detected efficiently. Also, radiation image read-out apparatuses using the long photomultiplier can be made smaller than apparatuses using a photomultiplier in which a light guide member having a complicated shape is located so that it is in close contact with a small light receiving face, as disclosed in, for example, Japanese Unexamined Patent Publication No. 54(1979)-87808.\nHowever, the aforesaid long photomultiplier has a drawback in that a large number of photoelectrons cannot readily be emanated from the photocathode. Specifically, in order to generate a large number of photoelectrons in the photocathode, it is necessary to increase the light absorption efficiency of the photocathode. For this purpose, the photocathode should be made thicker. However, photoelectrons generated at positions deep within the width of the photocathode cannot readily be emanated out of the photocathode. From this viewpoint, the photocathode should be made thinner. These two incompatible requirements make it difficult for a large number of photoelectrons to emanate from the photocathode.\nThe problems described above with regard to the detection of light emitted by a stimulable phosphor sheet by using a long photomultiplier also arise when other types of light are detected with the long photomultiplier."}
{"text": "When applications attempt to play more than one media stream on current devices, all the applications are allowed access to the presentation devices of the device, for example a display device and/or an audio device. The media streams are played by corresponding applications without regard for other media streams being played. Watching a video or listening to a song with interference from other audio streams and video streams is a common experience.\nWhen listening to a song and browsing the web, many web sites include audio in their web pages. The web page audio plays despite the fact that a song is already playing. This often leads to an unpleasant listening experience. If a user locates multiple videos and accesses them in multiple browser windows and/or tabs, the videos play as if a user is able to watch all of them at the same time. Videos in windows that are obscured by other windows or that are minimized continue to play as if there was someone watching. Some web pages do wait to detect they have input focus before beginning to play a stream, but these pages play their streams without regard for other media players playing and/or otherwise accessing a display or speakers to play one or more media streams.\nAccordingly, there exists a need for methods, systems, and computer program products for coordinating playing of media streams."}
{"text": "1. Field of the Invention\nThe present invention is generally directed to computer animation. Specifically, the present invention is directed to customizing a computer animation wireframe with three-dimensional range and color data or with a two-dimensional representation and a depth map.\n2. Introduction\nKnown systems can provide generic computer animations that integrate audio and visual information. For example, these generic computer animations typically display a talking head of a human or of a cartoon animal. These generic computer animations can be used for a number of applications.\nFor example, some known systems display the computer animation on a computer video monitor to interface with a human user. Other known systems can convert ASCII (American Standard Code for Information Interchange) text into synthetic speech and synchronized talking-head video with realistic lip and jaw movements.\nThese known computer animations are based on generic animation wireframe models. Although these generic animation wireframe models are generic in the sense that the animations doe not represent a specific person; these generic models can be deformed according to a predefined set of parameters to vary the presentation from the one generic version. Deforming a generic animation wireframe model can be used to more closely resemble realistic and natural interactions, for example, human-to-human interactions. Deforming the generic model using a predefined set of parameters, however, cannot sufficiently modify the generic model to present actual people recognized by the viewer.\nTo produce more realistic and natural displays for human interactions, animation wireframe models should incorporate real measurements of the structure of the desired face, as well as color, shape and size. Such information can be obtained by a three-dimensional laser scanner system that scan a person's head to produce very dense range data and color data of the head.\nSome known systems that incorporate measured three-dimensional information into generic animation wireframe models, however, suffer from several shortcomings. In general, accurately modifying generic animation wireframe models with measured three-dimensional range data requires extensive and expensive manual adjustments or automated computer-based adjustments. Manual adjustments of generic animation wireframe models can be time consuming and/or can require expensive human personal with specialized training. Automated adjustments of generic animation wireframe models can require expensive computer equipment that is generally cost-prohibitive for mass distribution and may require extensive maintenance performed by human personnel with specialized training."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an information processing apparatus of an on-vehicle type such as an audio apparatus of the on-vehicle type, and more particularly to an information processing apparatus of an on-vehicle type, which is mounted on a vehicle and which has a fan for cooling and/or removing dust.\n2. Description of the Related Art\nConventionally, it is general that an audio apparatus such as a so-called car-radio, a cassette deck, a CD (Compact Disc) player or the like is mounted in a room or cabin of the vehicle, so as to serve music etc., to the users or passengers on the vehicle by the audio apparatus.\nRecently, a so-called navigation apparatus is also generalized, which has a display device such as an LCD (Liquid Crystal Display) device or the like in the cabin of the vehicle and helps a movement of the vehicle by displaying a necessary map etc., on the display device.\nIn such an audio apparatus or a navigation apparatus, a device which generates heat during the operation thereof (e.g., a CPU, an amplifier or the like) and a device which should not be heated up to a certain high temperature during the operation thereof (e.g., an MD (Mini Disc) drive or the like) may be commonly built-in. Thus, a fan for cooling is often built-in so as to restrain the temperature increase by heat from the device generating the heat.\nWhen the supply of the electric power with respect to the audio apparatus or the navigation apparatus is started, the rotation of the fan is started at the same time so as to perform the cooling operation while the rotation speed of the fan (which is defined as a rotation number per unit time, hereinbelow) is constant.\nThere is recently such a tendency that the silence within the cabin is regarded as an important factor.\nHowever, even if the cabin is designed to attach importance onto the silence, since the audio apparatus or the navigation apparatus is mounted in the cabin of the vehicle, the silence cannot be preserved because of the blowing sound due to the rotation of the fan within the audio apparatus or the navigation apparatus, which is a problem.\nOn the other hand, if the rotation speed of the fan is reduced since too much importance is attached onto the silence within the cabin, the cooling function of the fan cannot be fulfilled, so that the audio apparatus or the navigation apparatus may be failed due to the heat, which is another problem.\nIt is therefore an object of the present invention to provide an information processing apparatus of an on-vehicle type, which can operate a fan for cooling or removing dust while preserving the silence in a cabin of the vehicle and at the same time preserving the rotation speed necessary for the fan, which is high enough to appropriately function.\nThe above object of the present invention can be achieved by a first information processing apparatus of an on-vehicle type provided with: a fan; a detecting device for detecting an operation status of a vehicle; and a controlling device, such as a CPU or the like, for controlling a rotation speed of the fan in accordance with the detected operation status.\nAccording to the first information processing apparatus of the present invention, since the rotation speed of the fan is controlled in accordance with the operation status of the vehicle, it is possible to preserve the silence in the cabin of the vehicle as the unnecessary blowing sound of the fan is restrained by decreasing the rotation speed in case that the priority is given to the silence, for example. It is also possible to preserve the necessary rotation speed of the fan to fulfill the cooling function, the dust sucking function or the like, by increasing the rotation speed in case that the priority is not given to the silence (but to the prevention against a failure of the information processing apparatus, for example).\nIn this manner, it is possible to preserve both of the silence in the cabin and the necessary rotation speed of the fan, depending upon the operation status of the vehicle.\nIn one aspect of the first information processing apparatus of the present invention, the detecting device comprises a silence degree detecting device, such as a voice microphone or the like, for detecting a degree of silence within a cabin of the vehicle, and the controlling device sets the rotation speed to a low speed or stops a rotation of the fan if the detected degree of silence is higher than a threshold degree of silence, which may be set in advance or may be changed in operation, and sets the rotation speed to a high speed, which is higher than the low speed, if the detected degree of silence is not higher than the threshold degree of silence.\nAccording to this aspect, since the fan is rotated at the low speed or the rotation of the fan is stopped in case that it is silent in the cabin, and since the fan is rotated at the high speed in case that it is not silent in the cabin, it is possible to preserve both of the silence in the cabin by restraining the blowing sound of the fan and the necessary rotation speed of the fan.\nIn another aspect of the first information processing apparatus of the present invention, the detecting device comprises a vehicle speed detecting device, such as a travel distance sensor or the like, for detecting a moving speed of the vehicle on the basis of a vehicle speed pulse signal outputted in response to the moving speed, and the controlling device sets the rotation speed to a low speed, which is lower than a reference rotation speed, or stops a rotation of the fan if the detected moving speed is not higher than a threshold moving speed, which may be set in advance or may be changed in operation.\nAccording to this aspect, since the fan is rotated at the low speed or the rotation of the fan is stopped in case that the vehicle is moving slowly, it is possible to preserve the silence in the cabin more surely.\nIn another aspect of the first information processing apparatus of the present invention, the detecting device comprises a vehicle speed detecting device, such as a travel distance sensor or the like, for detecting a moving speed of the vehicle on the basis of a vehicle speed pulse signal outputted in response to the moving speed, and the controlling device sets the rotation speed to a high speed, which is higher than a reference rotation speed, if the detected moving speed is higher than a threshold moving speed, which may be set in advance or may be changed in operation.\nAccording to this aspect, since the fan is rotated at the high speed in case that the vehicle is moving fast, it is possible to preserve the necessary rotation speed of the fan if it can be assumed that it is not necessary to preserve the silence in the cabin as the vehicle is moving fast.\nIn another aspect of the first information processing apparatus of the present invention, the detecting device comprises a vibration detecting device, such as an angular velocity sensor or the like, for detecting a vibration of the vehicle, and the controlling device sets the rotation speed to a low speed, which is lower than a reference rotation speed, or stops a rotation of the fan if the detected vibration indicates the operation status that the vehicle is not moving.\nAccording to this aspect, since the fan is rotated at the low speed or the rotation of the fan is stopped in case that the vehicle is not moving, it is possible to preserve the silence in the cabin more surely.\nIn another aspect of the first information processing apparatus of the present invention, the detecting device comprises a vibration detecting device, such as an angular velocity sensor or the like, for detecting a vibration of the vehicle, and the controlling device sets the rotation speed to a high speed, which is higher than a reference rotation speed, if the detected vibration indicates the operation status that the vehicle is vibrating by a certain magnitude.\nAccording to this aspect, since the fan is rotated at the high speed in case that the vehicle is vibrating, it is possible to preserve the necessary rotation speed of the fan if it can be assumed that it is not necessary to preserve the silence in the cabin as the vehicle is vibrating.\nIn another aspect of the first information processing apparatus of the present invention, the information processing apparatus further comprises a sound detecting unit, such as a voice microphone or the like, for detecting a sound in a cabin of the vehicle, and the controlling device sets the rotation speed to a high speed, which is higher than a reference rotation speed, if a sound volume of the detected sound is not smaller than a threshold volume, which may be set in advance or may be changed in operation, regardless of a detection result of the detecting device.\nAccording to this aspect, it is possible to preserve the necessary rotation speed of the fan in case that it is not actually silent in the cabin regardless of the operation status of the vehicle.\nIn another aspect of the first information processing apparatus of the present invention, the information processing apparatus further comprises a sound detecting unit, such as a voice microphone or the like, for detecting a sound in a cabin of the vehicle, and the controlling device sets the rotation speed to a low speed, which is lower than a reference rotation speed, or stops a rotation of the fan if a sound volume of the detected sound is not larger than a threshold volume, which may be set in advance or may be changed in operation, regardless of a detection result of the detecting device.\nAccording to this aspect, it is possible to restrain the generation of the blowing sound of the fan, which disturbs the silence in the cabin, in case that it is actually silent in the cabin regardless of the operation status of the vehicle.\nIn another aspect of the first information processing apparatus of the present invention, the information processing apparatus further comprises a temperature detecting unit, such as a thermistor or the like, for detecting a temperature within the information processing apparatus, and the controlling device sets the rotation speed to a high speed, which is higher than a reference rotation speed, if the detected temperature is not lower than a threshold temperature, which may be set in advance or may be changed in operation, regardless of a detection result of the detecting device.\nAccording to this aspect, it is possible to cool the information processing apparatus by rotating the fan at the high speed, in case that the temperature in the information processing apparatus is actually high, regardless of the operation status of the vehicle.\nIn another aspect of the first information processing apparatus of the present invention, the information processing apparatus further comprises a temperature detecting unit, such as a thermistor or the like, for detecting a temperature within the information processing apparatus, and the controlling device sets the rotation speed to a low speed, which is lower than a reference rotation speed, or stops a rotation of the fan if the detected temperature is not higher than a threshold temperature, which may be set in advance or may be changed in operation, regardless of a detection result of the detecting device.\nAccording to this aspect, it is possible to preserve the silence in the cabin, by restraining the unnecessary blowing sound of the fan, in case that the temperature in the information processing apparatus is actually low, regardless of the operation status of the vehicle.\nThe above object of the present invention can be also achieved by a second information processing apparatus of an on-vehicle type provided with: a fan; a sound volume detecting device, such as a noise microphone or the like, for detecting a sound volume of a sound generated in a cabin of a vehicle; and a controlling device, such as a CPU or the like, for controlling a rotation speed of the fan in accordance with the detected sound volume.\nAccording to the second information processing apparatus of the present invention, since the rotation speed of the fan is controlled in accordance with the sound volume of the sound in the cabin, it is possible to preserve the silence in the cabin of the vehicle as the unnecessary blowing sound of the fan is restrained by decreasing the rotation speed in case that the priority is given to the silence, for example. It is also possible to preserve the necessary rotation speed of the fan to fulfill the cooling function, the dust sucking function or the like, by increasing the rotation speed in case that the priority is not given to the silence (but to the prevention against a failure of the information processing apparatus, for example).\nIn this manner, it is possible to preserve both of the silence in the cabin and the necessary rotation speed of the fan, depending upon the operation status of the vehicle.\nIn one aspect of the second information processing apparatus of the present invention, the controlling device sets the rotation speed to a low speed, which is lower than a reference rotation speed, or stops a rotation of the fan if the detected sound volume is not larger than a threshold volume, which may be set in advance or may be changed in operation.\nAccording to this aspect, since the fan is rotated at the low speed or the rotation of the fan is stopped in case that it is silent in the cabin, it is possible to preserve the silence in the cabin by restraining the unnecessary blowing sound of the fan.\nIn another aspect of the second information processing apparatus of the present invention, the controlling device sets the rotation speed to a high speed, which is higher than a reference rotation speed, if the detected sound volume is larger than a threshold volume, which may be set in advance or may be changed in operation.\nAccording to this aspect, since the fan is rotated at the high speed in case that it is not silent in the cabin, it is possible to preserve the necessary rotation speed of the fan in case that it is not actually silent in the cabin.\nIn another aspect of the second information processing apparatus of the present invention, the information processing apparatus further comprises a temperature detecting unit, such as a thermistor or the like, for detecting a temperature within the information processing apparatus, and the controlling device sets the rotation speed to a high speed, which is higher than a reference rotation speed, if the detected temperature is not lower than a threshold temperature, which may be set in advance or may be changed in operation, regardless of a detection result of the sound volume detecting device.\nAccording to this aspect, it is possible to cool the information processing apparatus by rotating the fan at the high speed in case that the temperature in the information processing apparatus is actually high, regardless of the sound volume of the sound in the cabin.\nIn another aspect of the second information processing apparatus of the present invention, the information processing apparatus further comprises a temperature detecting unit, such as a thermistor or the like, for detecting a temperature within the information processing apparatus, and the controlling device sets the rotation speed to a low speed, which is lower than a reference rotation speed, or stops a rotation of the fan if the detected temperature is not higher than a threshold temperature, which may be set in advance or may be changed in operation, regardless of a detection result of the sound volume detecting device.\nAccording to this aspect, it is possible to preserve the silence in the cabin as the unnecessary blowing sound of the fan is restrained by rotating the fan at the low speed in case that the temperature in the information processing apparatus is actually low, regardless of the sound volume of the sound in the cabin.\nThe above object of the present invention can be also achieved by a third information processing apparatus of an on-vehicle type provided with: a fan; at least two of an operation status detecting device, such as a noise microphone or the like, for detecting an operation status of a vehicle, a sound detecting unit, such as a voice microphone or the like, for detecting a sound in a cabin of the vehicle and a temperature detecting unit, such as a thermistor or the like, for detecting a temperature within the information processing apparatus; and a controlling device, such as a CPU or the like, for setting a rotation speed of the fan to a high speed, which is higher than a reference rotation speed, if a detected value of at least one of the operation status detecting device, the sound detecting unit and the temperature detecting unit is higher than a threshold value, which is set in advance for each of the operation status detecting device, the sound detecting unit and the temperature detecting unit.\nAccording to the third information processing apparatus of the present invention, since the rotation speed of the fan is set to the high speed in case that the detected value of either one of the detecting device is not less than the threshold value thereof, i.e., in case that the generation of the blowing sound of the fan is admitted, it is possible to preserve the necessary rotation speed of the fan, as the occasion demands.\nIn another aspect of the first information processing apparatus of the present invention, in another aspect of the second information processing apparatus of the present invention or in one aspect of the third information processing apparatus of the present invention, the information processing apparatus reproduces information supplied from a plurality of kinds of sources.\nAccording to this aspect, in case that the information is reproduced in the cabin, it is possible to preserve both of the silence in the cabin and the necessary rotation speed of the fan, by controlling the rotation speed of the fan as the occasion demands.\nThe nature, utility, and further features of this invention will be more clearly apparent from the following detailed description with respect to preferred embodiments of the invention when read in conjunction with the accompanying drawings briefly described below."}
{"text": "This invention relates to a steering torque detecting apparatus, and more particularly to a steering torque detecting apparatus which is simplified in structure, readily assembled and exhibits high reliability.\nNow, a conventional steering torque detecting apparatus will be described with reference to FIGS. 1 to 3. FIG. 1 shows a conventional steering system for an outboard motor. A steering handle 101 is provided so as to be rotatably operated in any desired direction. The rotation of the steering handle 101 is transmitted through a steering wire 102, a slide member 105, a steering rod 107 and an oscillation lever 109 to an outboard motor 111.\nSteering force input to the steering wire 102, as shown in FIG. 2, is also transmitted through an interlocking element 113 and a transmission rod 115 to a torque sensor 117. The torque sensor 117 includes a sensor body 121 slidably arranged in a housing 119 as shown in FIG. 3. The sensor body 121 is normally held at a neutral position by means of a spring 123. The sensor body 121 is formed at a part thereof into a reduced diameter, resulting in providing a reduced diameter section 125, on which a strain sensor 127 is mounted. The sensor body 121 is also connected to the transmission rod 115. To the strain sensor 127 is connected a signal cable 129, which is arranged so as to externally extend through a through-hole 131 formed at the housing 119.\nWhen the sensor body 121 which is connected to the transmission rod 115 as described above is forced in the longitudinal direction thereof against the spring 123, strain occurs on the reduced diameter section 125 of the sensor body 121. The so-produced strain is detected by the strain sensor 127 and then input to a controller 133 shown in FIG. 1. The controller 133 serves to calculate power assisting force depending upon a detection signal supplied from the strain sensor 127 to generate a control signal, which is then supplied to a drive motor 135, leading to rotation of the drive motor 135. The rotation of the drive motor 135 causes a pinion 137 (FIG. 2) to be rotated, so that a rack 139 is moved in a suitable direction.\nSuch movement of the rack 139 results in the power assisting force being applied through the steering rod 107 to the oscillation lever 109.\nUnfortunately, it was found that the conventional steering torque detecting apparatus constructed as described above has the following problems.\nFirst, the operation of mounting the strain sensor 127 on the torque sensor 117 is highly troublesome, because it must be carried out in a narrow space in the housing 119.\nAnother problem of the conventional steering torque detecting apparatus is that the strain sensor 127 is readily subject to a radio trouble or fault. This adversely decreases the level of output of the strain sensor to a degree sufficient to cause the apparatus to require an amplifier, resulting in the apparatus being highly large-sized and costly."}
{"text": "Within the United States of America and many other areas of the world, one of the most popular game balls is that known as an American football. The ball used in the American version of football is similar to the earlier ball used in the European game of Rugby but has evolved to a smaller size having generally more pointed end portions. Through tradition and a number of competitive league rules, footballs are generally fabricated in accordance with rigidly and precisely defined physical construction. While some variety of footballs exist in different footballs leagues, most include a spheroidal outer skin formed of leather or plastic within which an inflatable rubber bladder is supported. The bladder is inserted into the outer skin through an elongated slot which is thereafter closed with a plurality of laces similar to a shoelace. In later developments, footballs have been formed having a continuous outer rubber skin manufactured through the use of molding processes. When so made, however, most such footballs retain the outer characteristic and appearance of a more conventional traditional football including the presence of simulated raised laces. The laces on the football have become important in that they provide a portion of the gripping mechanism used by a football passer in throwing a football.\nIn attempts to find more interesting and exciting types of footballs for use by players of all ages, practitioners in the art have produced a variety of football designs which are \"nonstandard\" and which depart from the above-mentioned traditional shapes and thus are used more for amusement and entertainment outside of organized traditional football competition. Among the objectives in providing such a variety of football shapes has been the need for different flight characteristics, different gripping patterns, as well as the general ever present desire for variety among consumers. In attempting to meet and satisfy this need for variety, practitioners have provided a virtually endless array of football shapes and configurations.\nFor example, U.S. Pat. No. 4,887,814 and U.S. Pat. No. Des. 294,844 both issued to Winter set forth a GAME BALL for sport and recreational activities having a football-shaped body formed of a resilient elastically deformable material having channels associated therewith. The channels are generally helically wound about the ball from one end to the other and are intended to receive the user's individual fingers within the channel when the ball is gripped.\nU.S. Pat. Re. 33,449 issued to Martin sets forth a HELICALLY GROOVED FOAM FOOTBALL in which an elastic foam football defines an outer surface within which lengthwise spiral grooves having increasing width and depth toward the middle portion of the ball are formed for improved handling. The helical pattern of the grooves is selected to permit the user to insert the user's finger end portions into a single groove when the football is gripped.\nU.S. Pat. No. 2,194,674 issued to Riddel sets forth a FOOTBALL in which a conventional shaped football is formed of a plurality of rigid outer segments joined at a plurality of helical seams, each joint seam forming upwardly extending ridges upon the football surface. An inflatable bladder is received within the football and inflated to pressurize the ball.\nU.S. Pat. No. Des. 235,794 issued to Kroener sets forth a FOOTBALL WITH SPIRAL SEAMS in which a conventionally shaped football defines a plurality of spiral seams extending from one end of the ball to the other. The seams form inwardly extending grooves.\nU.S. Pat. No. 2,859,040 issued to Gow, et al. sets forth a FOOTBALL HAVING A SECURELY GRIPABLE LACELESS SURFACE in which a football defines a pressurized interior cavity having a pneumatic bladder therein and an outer surface characterized by a plurality of gripping ridges. The center one-third of the football defines a plurality of straight ridges extending along the football's major axis while the remaining end portions of the football surface define concentric circular ridges extending transversely to the football major axis.\nU.S. Pat. No. 4,736,948 issued to Thomas sets forth a FOOTBALL having an oblate spheroid body defining a passageway along its longitudinal axis. A pair of wind fins are mounted internally of the body so as to protrude into the passageway. An alternate embodiment is set forth in which a plurality of inwardly extending grooves are defined in the outer surface of the football.\nU.S. Pat. No. 3,884,466 issued to MacDonald, et al. sets forth a GAME BALL having a football shape which defines an air passage extending through its longitudinal axis. The diameter of the air passage tapers from opposite ends of the football to a constricted opening in the center thereof. A plurality of relatively heavy weights are supported by the football encircling the constricted passage.\nU.S. Pat. No. 1,931,429 issued to Buckner, et al. sets forth a FOOTBALL having a conventional inflated football further defining a plurality of helical grooves in the outer surface thereof. An abrasive material is deposited within the helical grooves to enhance the gripping characteristics of the football.\nU.S. Pat. No. 4,919,422 issued to Ma sets forth a CURVE BALL having a generally spherical shape and balance and defining an axis therethrough. First and second convex surfaces are defined about the axis at opposite poles thereof and at least one groove having a non-uniform depth defined in the surface of the ball between the first and second convex surfaces. The groove is arranged to extend through the equatorial region of the ball between the first and second convex surface. The spherical-shaped ball is thrown so as to impart spin thereto at different angles to the axis to create various curving actions.\nWhile the foregoing described prior art devices have provided increased variety and different aerodynamic characteristics for football shaped game balls, there remains a continuing need in the art for evermore interesting and varied game balls to meet consumer appetites for improvement and variety."}
{"text": "This invention relates to an automatic slack adjuster of the type used for commercial vehicle braking systems.\nCommercial vehicle braking systems typically incorporate a slack adjuster in each brake assembly to adjust the clearance between the brake linings and the rotating brake element, such as a brake drum or rotor. The brake assembly includes an actuator, such as an air chamber, that urges the brake linings into engagement with the brake element. As the linings wear, the clearance between the brake linings and brake element increases requiring the air chamber push rod to move a greater distance to apply the brakes. It is desirable to maintain a relatively constant clearance throughout the life of the brake linings to provide consistent braking performance. Slack adjusters are employed to keep the distance that the air chamber push rod must move within a specified range as the linings wear to maintain a consistent clearance.\nAutomatic slack adjusters have been developed to adjust for the clearance between the brake linings and brake element during normal vehicle operation. One automatic slack adjuster available from the Assignee of the present invention incorporates a worm gear in engagement with a gear. The gear is connected to a camshaft that moves the brake lining into engagement with the brake element upon actuation of the air chamber. When excess clearance occurs, the worm gear is rotated by an adjustment assembly connected between the worm gear and the push rod. Rotating the worm gear adjusts the rotational position of the slack adjuster relative to the camshaft, which adjusts the push rod travel.\nThe worm gear is disposed within a bore in the housing. An end portion of the worm gear has longitudinally extending teeth. A cylinder-shaped actuator is disposed within the bore and includes an inner diameter having longitudinally extending teeth engaging the teeth of the worm. An actuator piston is disposed within an internal actuator cavity and is retained therein by a piston retaining ring. An end of the actuator rod is pinned to the actuator piston. A spring loaded pawl assembly is supported by the slack adjuster housing and includes an end having teeth that engage the outer diameter of the actuator. The outer diameter of the actuator has helical teeth that cooperate with the teeth on the pawl assembly.\nThe actuator rod moves the piston along a length defined by the actuator cavity in response to the brake being applied and released. The actuator fits loosely within the bore to permit lateral movement of the actuator within the bore. When excess clearance has developed, the actuator rod will pull the actuator in a direction away from the worm gear with the actuator piston. As a result, the actuator will “jump” a tooth relative to the pawl assembly teeth such that on the brake release the worm will rotate the gear relative to the slack adjuster for the next brake apply thereby taking up the clearance.\nAs can be appreciated for the above description of prior art slack adjusters, the slack adjuster utilizes numerous components adding cost and complexity to the assembly of the slack adjuster. For example, the actuator requires machining on both the inner and outer diameters, and the spring loaded pawl requires numerous parts. Therefore, what is needed is a simplified slack adjuster that reduces the cost of the assembly."}
{"text": "This invention relates generally to computed tomography (CT) imaging and more particularly, to quantification of a selected attribute of an image volume and monitoring changes of the selected attribute in a patient.\nVisualization of anatomical data acquired by imaging devices generating 3D data is typically handled by volume rendering of its intensity and/or density values (Hounsfield Units (HU) in the case of Computed Tomography for instance). Many clinical applications are based on 3D visualization of the volumetric data that may include, but are not limited to, detection and sizing of lung nodules, quantification of vessel curvature, diameter and tone, cardiac vascular and function applications, and navigation of the colon for detection of polyps. These applications rely on the absolute values of the image data (intensity, density (HU), uptake (standard uptake values (SUV)), and other material properties associated with medical imaging to differentiate multi-dimensional anatomies from background tissue. Some clinical imaging applications are designed for routine screening of early cancers in the form of, for example, tumors, nodules, and polyps.\nMany cancers commonly metastasize or move from their primary organ or location to involve another organ or location. The most common location for tumors to metastasize to is lymph nodes followed by lung, liver, and then bone. Frequently, metastatic disease presents as a distribution of small lesions (2-10 mm) throughout the anatomy of the body. Most common locations for metastatic lesions are in the lung and liver. The visual contrast of liver lesions on CT images is limiting to the human eye. Magnetic Resonance Imaging (MRI) and Positron Emission Tomography (PET) imaging have proven superior to CT for visualizing liver tumors, but contrast remains limited.\nThere are many treatment options for primary and secondary cancers. These may include radiation therapy, chemotherapy, hormone treatment, immune therapy, surgery and others.\nTo date, physicians rely significantly on the apparent anatomical size and shape of the tumor under treatment when assessing the patients response to a chosen therapy. This can be problematic in patients with “bulky disease” (meaning that the tumor burden is an overestimate of the actual presence of cancer cells) if the cancer remises, but the relative size of the tissue mass does not change. Since the inception of PET and CT/PET imaging, the size of the active portion of the tumor can be when assessed to determine patient response to therapy. The physician may desire to measure the size of the lesion(s) before and after subsequent treatments to quantify the response. In many cases of a primary cancer, it may be straightforward to quantify the volume of anatomy occupied by a lesion. Under some circumstances, a tumor can have limited contrast and/or be ill-defined, meaning that the boundaries of the tumor are difficult to identify. In the case of multiple lesions and metastatic disease, there may be hundreds of small lesions distributed throughout the body or within individual organs. However, when there are multiple lesions, it is extremely time consuming to identify and track each individual lesion. Additionally, physicians may choose to represent the sum total of the volume occupied by all of the lesions in terms of a single number called “Total Tumor Burden” (TTB). As such, when any of the tumors respond to a chosen treatment plan, the TTB will change. However, even tracking TTB over the course of a treatment regime may also require a difficult and time consuming procedure."}
{"text": "In the fields of molecular and cellular biology, particularly in a laboratory setting, it is of vital importance that the correct reagents and media are used in every instance, including use on the bench or in a sterile hood, and when selecting, ordering or re-ordering supplies. There are various factors that exist in a laboratory setting that can make this more difficult and less efficient. For example, bottles of media and sera are often stored in refrigerators where condensation may cause labels to be hard to read, and the important information may be small print that is easily obscured by moisture, user writing to label and identify additions, or other information about a particular bottle such as the user's initials, date opened, etc. Another difficulty is the manual dexterity required for sterile transfer of materials from a reagent or media bottle. Laboratory workers routinely wear protective gloves and must hold a pipette in one hand while also holding and opening a bottle, all while maintaining the sterility of the pipette, the cover and the contents of the bottle. Prior to the present disclosure, it has been difficult to quickly identify with assurance the contents of a particular bottle. This concern is of particular importance in the context of cell culture media bottles where mixing culture media types or solutions can have disastrous effects on research projects, greatly increasing the cost and time associated. Often, different types of media are indistinguishable by merely looking at the contents, requiring reading of small print, or re-orienting a bottle in order to positively identify the contents.\nAlthough conventional labels often include the name of the product family in large type on the front of the bottle, it may be difficult to distinguish exactly what product is in the bottle, particularly if there are user annotations covering part of the label. For example, a certain label may identify the contents as a certain media, such as a basal media, but it may not be as clear that the content is basal media+glucose. Extra time and extra steps are needed, therefore, to determine exactly which product is being used, and to find a catalog number for re-ordering. Although a particular instance of wasted time might be small, they can collectively consume excessive amounts of time when one considers that within a single laboratory these actions may take place tens or even hundreds of times within a day.\nIn addition to difficulty in identifying the contents of a bottle, conventional bottles typically do not offer both storage efficiency and versatility. As all storage devices, refrigerators, incubators, freezers, shakers, stirrers, hoods, culture rooms, etc., do not have the same layout and dimensions, adaptability and versatility in the storage capabilities of bottles can be a significant problem. One type of conventional bottle is the round bottle shown in FIG. 1. This bottle shape results in inefficient use of space when stored, due to dead space between bottles and the inability to stably stack bottles on top of each other. Another type of conventional media bottle is a square bottle as shown in FIG. 2. Although these bottles can more efficiently abut other bottles both from the side and the front, they cannot be stored on a side and stacked on top of one another without the contents of the bottle coming in contact with the lid. Content contact with the lid can potentially cause leakage or contamination resulting in increased costs, lost research time, and a lack of supplies. Another disadvantage is that when the bottles are stored side by side in rows, only the labels of the front row can be seen without lifting the bottles to reveal a label. Furthermore, the caps are not easily accessible and the tops of the caps cannot be seen without lifting the bottle or looking directly down from above, which is not possible in many storage situations such as in a refrigerator. As all storage devices, refrigerators, incubators, freezers, shakers, stirrers, hoods, culture rooms, etc., do not have the same layout and dimensions, adaptability and versatility in the storage capabilities of bottles can be a major source of inefficiency.\nPrior to this disclosure, there have been few or no attempts to solve these problems, especially in the fields of microbiology and cellular biology where conventional round and square media bottles are almost universally used. A mixing station for chemicals used in cleaning and maintenance supplies using a color coded system has been disclosed in U.S. Pat. No. 6,322,242. This patent discloses the use of various color coded elements, including caps, labels, valve members, and liquid outlet lines, to identify how the components should be assembled with respect to one another. This does not teach, however, a system of providing solid or liquid products that are appropriate or useful for microbiological or other research media, particularly to be used with sterile technique, and utilizing color coded covers or caps, corresponding color coded labels, or specially designed bottles to facilitate ease of identification of the contents of a bottle during storage, use, and re-supplying."}
{"text": "The evolution of saddles has recently started to gather pace. From the very early saddles, designed purely to provide a more secure seat for the rider, through the development of stirrups, leaving aside pack saddles, there have been three main strands of design: the military, designed to secure the rider firmly and provide a degree of protection; the working or western saddle, also designed to provide a degree of security for the rider; and the English or close contact style where security for the rider is subservient to need for close contact with the animal. All three styles have a solid backbone or “tree”, traditionally made of wood (more recently materials have included fiberglass, metal and plastic), round which the leather (or synthetic equivalent) is mounted. Inevitably such a rigid structure placed on a moving surface raises difficulties with the fit of the saddle to the horse. In the case of the military and western saddles this is partially addressed by using a thick saddle blanket. However with close contact saddles the issue is addressed by attention to fit, either by having a bespoke saddle made for the horse (which is very expensive), or by careful selection from a range of off the peg designs. It is estimated that a saddlery wishing to carry a basic range of off the peg saddles, covering the three main saddle styles (dressage, jumping, general purpose), in one single colour option, and to fit most sizes of horse and rider, would have to stock in excess of 72 different saddles.\nEven when a rider invests in a bespoke saddle, the traditional, static design based on a rigid tree does not allow for the changes in a horse's shape that occur as it moves, or as there are variations in its fitness. Even the best fitting saddle cannot distribute the pressure evenly throughout the range of a horse's movement, and even a well fitting treed saddle will inevitably create pressure points on the horse's back, especially when turning tightly, where the saddle tree acts somewhat as a splint longitudinally on the spine, or when riding up or down hill or jumping, where the load is focused by the tree towards the front or back of the saddle. This can cause pain and restrict movement, and can potentially leading to a range of physiological and behavioural problems such as bucking, rearing, lameness, bruising of the muscles, muscular atrophy and in more severe cases, tissue necrosis.\nOver the last thirty years several new designs of saddle have been developed, both to try to address the problems enumerated above, and to facilitate newly evolved riding disciplines such as endurance and vaulting. All still use a static method of mounting the saddle on the horse. Many of these new designs are described as “treeless”, but in practice most are semi-treed, in that they have a rigid internal fitting at either the pommel or the cantle of the saddle. This can lead to weight being distributed over fewer points than a standard tree, which, in some circumstances, can exacerbate the problem. Saddles that have no tree at all do nothing to spread the pressure of the girth and the stirrups, the full force of which is therefore concentrated immediately over the mounting points. There is also a perception that such saddles are not as secure on the horse, as many treeless designs do not include a gullet, which has the effect of reducing lateral stability. A further disadvantage of many such saddles is that it is difficult to design them to look like the traditional English saddle, a look that is very popular in the market.\nAn additional issue with traditional close contact saddle design is that the mounting position of the stirrups can be quite critical to the ability of the rider to effectively balance on their horse. Many buyers' choice of saddle is primarily based on this factor, in an attempt to ensure that they are able to sit in the ideal position “over” the stirrups. There is little or no allowance in most saddles for any adjustment of the stirrup bar mounting position, so that this factor can quite severely restrict the choice of saddle, and associated ability to ensure a good fit.\nIn the following discussion, the invention will be generally described in relation to equestrian uses of the invention. However, the invention is broadly applicable to pack animals as well as mounts for personal transport.\nAn object of the invention is to obviate or mitigate at least some of the aforesaid problems by providing improvements in saddle design."}
{"text": "1. Field of the Invention\nThis invention relates to methods and an apparatus for producing decabromodiphenyl alkanes. More specifically, the field of the invention is that of producing decabromodiphenyl ethane.\n2. Description of the Related Art\nHalogenated aromatic compounds are often employed as flame retardant agents. Flame retardants are substances applied to or incorporated into a combustible material to reduce or eliminate its tendency to ignite when exposed to a low-energy flame, e.g., a match or a cigarette. The incorporation of flame retardants into the manufacture of electronic equipment, upholstered furniture, construction materials, textiles and numerous other products is well known.\nBrominated aromatic compounds are often utilized as flame retardant agents in polymer compositions such as the outer housing of computers, television sets, and other electronic appliances. One group of halogenated flame retardants are decabromodiphenyl alkanes. The manufacture of decabromodiphenyl alkanes is known. Conventionally, decabromodiphenyl alkanes are prepared by reacting a diphenyl alkane with bromine in the presence of a bromination catalyst, such as AlCl3 or FeCl3.\nFor example, U.S. Pat. No. 5,030,778 to Ransford discloses a process for producing decabromodiphenyl alkanes in which bromine and a bromination catalyst are charged to a reaction vessel. Liquid diphenyl alkane is fed by a dip tube into the reaction vessel at a point which is beneath the level of the charged liquid bromine and catalyst. The stated advantages of this sub-surface addition method are that (1) a product with a high average bromine number is obtained faster when the diphenyl alkane is fed below the surface of the charged liquid bromine and catalyst; and (2) splattering of the reaction mass associated with the addition of the diphenyl alkane into the vessel is reduced.\nOne disadvantage to adding the diphenyl alkane to the vessel at a location below the surface of the charged bromine and catalyst is that the dip tubes used for adding the diphenyl alkane to the vessel are prone to plugging. It is believed that the sub-surface addition dip tubes become plugged when a small amount of diphenyl alkane remains at the tip of the tube and reacts in place, thereby forming insoluble, high melting point material. It is believed that this is more likely to occur at the end of the addition or if the diphenyl alkane addition is interrupted. This susceptibility to plugging prevents the manufacturer from being able to stop and start the diphenyl alkane addition, which is sometimes desirable for controlling the evolution HBr gas.\nIt is also believed that the agitation of the reaction mass may create a vortex within the tip of the sub-surface addition dip tube. This vortex may pull solids from the reaction mass into the tube, thereby creating a blockage. Additionally, because some diphenyl alkanes, such diphenylethane (“DPE”), are solids at room temperature and are fed to the reaction vessel as liquids, they may begin to crystallize in an unheated dip tube if the feed is interrupted for any reason. In the event that the sub-surface dip tube does become plugged, regardless of the reason, the diphenyl alkane feed must be stopped and the tube must be pulled out of the reactor in order to remove the blockage. It is desirable to avoid the need to remove the dip tube, as the vapor space of the reaction vessel is filled with toxic and corrosive bromine vapors which may escape during removal.\nThe above-surface diphenyl alkane addition technique of the present invention reduces dip tube plugging, thereby providing a more efficient method of adding diphenyl alkane to a reactor charged with bromine and catalyst."}
{"text": "In conventional chair structure, the chair cushion is usually parallel with the floor, and the angle between the cushion and a backrest is usually less than 90 degrees, so the user may easily feel physically/physiologically uncomfortable, and the user may not sit on the chair for a long period of time. Furthermore, the user may suffer from pain or sickness if sitting on the chair on a regular basis. Therefore, there remains a need for a new and improved chair structure to overcome the problems stated above."}
{"text": "Entrepreneurs engaged in a variety of business endeavors have found considerable utility in the use of motor vehicles generally referred to as \"pick-up\" trucks. These trucks are relatively light as compared to typical heavy haulage vehicles, having capacities ranging from about 1/2 to 3/4 ton. To enhance the appeal of the pick-up trucks, manufacturers have developed pleasing, streamlined body styles, for example, the lines of the bed portion of the trucks generally are integrated in flowing fashion with the lines and design of the cab and hood structure and ornamentation thereupon.\nFor many uses, the pick-up truck operators have developed a need for power dumping capability. However, while the beds of the trucks are readily removable by unbolting the floor thereof from the truck undercarriage, dump hoist installations heretofore have required the remounting of the dump bed at a higher level above the undercarriage in order to provide hoist access space. For example, the differential gear housing often is found to interfere with hoist installations. In addition to the readily apparent degredation to the design lines and aspects of the vehicle, such elevation of the bed derogates from vehicle stability under load. Vehicle safety further is compromised, inasmuch as the bed is only coupled to the truck body at rearward disposed hinges. In the event of a rearend accident or the like, the hinges tend to be torn from the body frame and the bed is thrust toward the cab. Degredation to body styling additionally is encountered in the repositioning of the rear bumper of the trucks following conversion. As is apparent, the dump bed cannot clear a bumper mounted to the termini of the frame side rail. Consequently, various schemes have been employed for remounting the bumper upon the frame at a lower level to provide dump clearance. Since bumpers are styled and contribute to the overall pleasing appearance of the trucks, this repositioning detracts from the looks of the converted truck.\nConversion heretofore carried out additionally have encountered problems relating to the structural integrity of the hoist equipment. In view of the relatively smaller size of the trucks, the hoist equipment itself remains small and the price of its installation must be low enough as to be commensurate with the lower cost of the trucks themselves. In consequence, the mere scaling down of conventional hoist structures to meet the needs of a smaller vehicle generally has been found to evolve structurally inadequate hoists. For example, extensive stressing has been determined to be present at the coupling between hydraulic motor piston rods and hoist components."}
{"text": "1. Field of the Invention\nThis invention pertains particularly to a drawing mixture which is suitable for rigorous utilizations involving maximum reductions and pass speeds, and also to an improved drawing procedure and utilization of drawing lubricants. A phase of the invention deals with a dry lubricating type of drawing mixture that has a proportioned content in a mixed relation and whose ingredients serve in a complementary highly improved manner to protect the surface of a metal workpiece being drawn under strenuous conditions.\n2. Description of the Prior Art\nVarious types of dry and wet lubricating compounds and mixtures have been used in the drawing of metal workpieces, such as steel, and in drawing high carbon and specialty steel material into wire or rod lengths. There has been a somewhat general adoption of a method which involves the application of a dip-applied lead coating to the workpiece to which is applied ordinary soap powder at the entrance to each draw die. In this connection, a cleaned rod, for example, is dipped to provide the lead coating and is then passed through a requisite number of dies to produce the desired final or semi-final product. In utilizing such a coating, it is important to control the speed so as to prevent a rise of temperature to near the melting point of the coating which in the case of lead is around 621.degree. F. Also, as the drawing progresses, the coating tends to become thinner and less effective thus limiting the percentage of reduction during continuation of the operation. For a less stringent type of operation, where the workpiece is not quite so hard or brittle a material, a modified form of coating is provided which utilizes an acid salt and may include a drawing compound in the nature of a sodium resin silicate.\nIt has been my experience that the present types of drawing or lubricating methods have all been of a production limiting nature, and although rejects have been reduced from about 50% to possibly 30% at the present time, the percentage is so great as to make the resultant product relatively expensive. Also, restrictions on the speed of operation, as well as the percentage reduction, have been limitations which have increased the cost of a product and made it relatively expensive. Also, restrictions on the speed of operation, as well as the percentage reduction, have been limitations which have increased the cost of a product such as welding rod, rope and bridge wire, music spring and valve spring wire, binding wire, etc."}
{"text": "Locks are used in distributed software to control access to shared resources. In some cases, a distributed lock manager (DLM) provides distributed applications with a mechanism for synchronizing their accesses to shared resources. For example, DLMs are used in some clustered file systems for file locking and for coordination of other disk accesses. The resources that are accessed under locks granted by a DLM may be files, records, areas of shared memory, peripherals, databases or portions of a database, and/or other items that are shared (or potentially shared) by different software processes.\nA lock may provide different kinds of access to a given resource, depending on the accesses already granted and the operations desired, for example. Different kinds of access are sometimes indicated as lock modes. For example, a lock request specifying a concurrent read mode could be used to indicate a desire to read a resource identified in the request without updating (writing) the resource. As another example, an exclusive mode indicates a desired for exclusive access which allows both reads and updates to the locked resource and prevents other software processes from having any access to the resource until the exclusive lock is released. Other lock modes may also be available in a given implementation of a DLM."}
{"text": "Scanning machines are often used to scan objects for the purpose of inspecting the scanned objects. Known scanning machines, such as the one disclosed in U.S. Pat. No. 4,020,346, include a conveyor to lead an object to be inspected into a scanning chamber where the object is subjected to radiation. A lead curtain is provided at an entrance and at an exit of the scanning chamber to contain the radiation within the chamber while allowing the passage of the object when it is conveyed in or out of the chamber.\nKnown to the Applicant are United States patent applications having Publication Nos. 2007/0133742, 2008/0025470, 2013/0114788 and 2013/0336447. Also known to the Applicant is U.S. Pat. No. 7,706,507 and Japan Patent No. 3,946,612. Each of the mentioned patents and patent applications disclose systems for containing radiation within the scanning chamber while allowing the passage of an object when it is conveyed in or out of the chamber.\nMoreover, a known process for the detection of explosive material involves manually taking a sample of fine particles from an object in order to subject the sample to a detection process. Typically such fine particles are collected with a cloth, which is then analyzed by an analyzing system which may include, for example, a heating chamber where the temperature is raised to gasify the particles containing explosive chemicals, from which the resulting gases are then analyzed with a spectrometer in order to detect the presence of such explosive chemicals.\nHowever, the teachings of the aforementioned suffer from drawbacks. For example, the inspection of an object involves the separate processes of scanning the object and detecting explosives, which can be both complicated and time consuming. Preferably, these processes should be simplified in order to make the inspection of objects more efficient.\nHence, in light of the aforementioned, there is a need for an improved system which, by virtue of its design and components, would be able to overcome or at least minimize some of the drawbacks of the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to the Vello process and draw-down process for making glass products from a glass melt and, more particularly, to a sleeve-type agitator for the glass melt for preventing schlieren formation and for homogenizing the glass melt in the Vello process and draw-down process. The invention also relates to a method and apparatus for feeding a glass melt in order to make glass products, especially glass tubes or pipes, from the glass melt without schlieren formation in the glass melt.\n2. Description of the Related Art\nThe Vello tube drawing process is a vertical drawing process for glass tube or pipe, in which the tube or pipe is drawn down from a circular nozzle. The glass melt flows through a nozzle with a cylindrical opening which is located at the bottom of a feed channel and over a funnel-shaped conical body which is adjustable in its height and which widens downward. The conical body is hollow and connected with a longitudinally extended pipe to a blower for air. The glass mass located around this valve body is held open when air is supplied through this pipe. The freely suspended glass is then drawn downward as a glass tube into a temperature regulated compartment or guided with the help of guide members into a horizontal orientation and drawn from the drawing machine.\nThe draw-down process operates exactly the same way that the Vello process does, except that the glass tube is not guided into a horizontal orientation, but instead is removed vertically from the drawing machine.\nIn the Vello head and draw-down head used in these glass product manufacturing processes schlieren are produced on the side facing away from the channel, which currently can be reduced only by connecting a head overflow. It is necessary to draw off the schlieren extending around the product with a head overflow duct in conventional feed heads. Faulty products would otherwise be produced with the schlieren extending around the product. Drawing the product with the schlieren thus causes glass losses."}
{"text": "1. Field of the Invention\nThis invention relates to a method of transferring data between computer systems with specific application in teleconferencing software programs. More particularly, it involves a method for transferring large amounts of data among interconnected computer systems according to the designated priority of the data, and for allowing the priority of the data to be changed before the data is completely transferred to a remote computer and for causing the remaining data to be subsequently transferred according to its new priority.\n2. Brief Description of the Prior Art\nWhenever two or more people are involved in the preparation of a document, whether it be a financial spread sheet, a CAD design, a circuit schematic layout, an organization report, a bit map image, etc., succeeding drafts of the document are prepared, circulated, and modified in the process. Each person annotates his or her remarks on the document and forwards it to the next person. Typically, several drafts of the document will be circulated before a final draft is produced. This is a very time consuming process.\nIn the case where a person involved in the document preparation process is at a different geographical location, getting the document from one location to another location and back becomes another tedious and time-consuming task. The document will either have to be mailed or faxed to that person, further complicating the entire process.\nOne standard method of alleviating this process is to hold meetings where everyone gathers and comments on the document with the hope of reducing the number of drafts needed before a final draft is produced. The shortcoming of this method is that there may be significant travel time and travel cost in getting all of the people to the same location. In addition, the final draft of the document is usually again circulated for final comments.\nOne solution to this problem is to use a teleconferencing software program, an aspect of which may contain an embodiment of the present invention. By using computer network connections or modem connected phone lines, everyone can be connected via his or her computer. By using the teleconferencing software program, everyone's computer screen displays the same data. In addition to using the software program and network or modem connections, conference calling over the voice phone lines or through the software program creates a dynamic and live atmosphere where everyone can participate in the discussion and refer to the document displayed on the screen.\nA very important capability of such teleconferencing software must be the ability to allow transfer of data from one computer user to other computer users. For example, in making a presentation using a number of frames of prepared graphs, charts, outlines, etc., each frame of data that is used must be quickly transferred to other users in the conference in order to have a common reference point for simultaneous discussion of the data presented. In addition, the presenter in the presentation may wish to skip among several frames of data or skip a few frames of data entirely. The teleconferencing software must allow this type of flexibility and still maintain a high efficiency in transferring data. At the same time, each frame of data must be organized in a manner that allows quick access by the users.\nAdditionally, the presenter may wish to transmit a private message to one particular user. The teleconferencing software will have to be able to distinguish between data for all interconnected computer systems (public data), and data for a particular user (private data), and properly transfer the data to the designated user or users.\nAnother problem in teleconferencing software is that the presenter may scroll through frames of data without allowing adequate time for the data to be transferred to all the computer systems. The presenter will eventually display one frame of data for discussion. At this time, this frame of data has the highest priority and must be immediately transferred to all other computer systems even though there may be several preceding frames of data that have not been completely transferred.\nNevertheless, all public frames of data scrolled through or loaded for the presentation must be transferred to all the connected users, because the presenter may eventually go back to previous frames of data in making his or her presentation. Thus, it is important to have the ability to organize the different frames of data and transfer the frames of data currently being used first, while establishing a system where other frames of data can be prioritized and transferred.\nAnother problem the present invention must deal with is the problem of transferring data between interconnected computer systems. The computer systems may be interconnected via modem, network, cellular links, or any other available connections. In connecting computer systems or nodes to computer systems, one computer may not be directly connected to all other computers involved in the teleconference. For example, referring to FIG. 1, there are four computer systems involved in this session of teleconferencing: computer A (10), computer B (12), computer C (14), and computer D (16). Computer A is only connected to computer B, computer B is connected to computer A and computer C, computer C is connected to computer B and computer D, and computer D is only connected to computer C. Computer A is connected to computer D only through computer B and computer C. In order for computer A to transfer data to computer D, the data must pass through computer B and computer C. Thus, if the user at computer A is making a presentation by using frames of data, these frames of data must travel through computer B and computer C to reach computer D. However, the data has to be transferred in such a manner so that there will not be a significant time lag between the time computer B receives the data, computer C receives the data, and computer D receives the data, so that all the users can follow the discussion or presentation in real-time or near real-time."}
{"text": "Frequently, homes, offices and industrial plant facilities experience many types of emergency situations involving power failures where an interior or exterior area is rendered without light. Such power failures may result from electrical short circuits, brownouts, fire, accidents, natural disasters (e.g., floods, hurricanes, tornados, etc.) or a planned shutdown of electricity within a facility or dwelling. Automotive commuters also frequently find themselves without power following accidents, vehicle break downs, and the like.\nIn each of the circumstances above, it would be desirable for person to possess a portable light device that is adapted to provide a source of local illumination and electrical power. However, the preferred type or mode of illumination may change depending upon the specific power outage circumstance. For example, it is generally preferable to have a broadly ranging flood-light type of illumination to reveal a person's path as they attempt to transit a darkened room or corridor. Alternatively, it is generally preferable to have a compact spot-light type of illumination to reveal a person's work area as they attempt to fix a flat-tire along a darkened road side.\nIn view of the above, a need exists to provide a portable light device that is capable of activation in response to a disruption of power. It would be desirable for the device to be compact and lightweight such that it can be moved simply and quickly from location to location. It is further desirable for the device to be efficiently powered by a battery pack such that it is capable of constant illumination lasting for several days or even weeks. It is further desirable that the portable light device be readily adapted to provide a spot-light mode of illumination, a flood-light mode of illumination, or a combination thereof. Finally, the device should provide real-time battery life information to a user such that the performance of the device may be tailored to extend or shorten expected battery life as needed."}
{"text": "1. Technical Field\nThe present disclosure relates to a light-emitting diode (LED) manufacturing method, and particularly to a manufacturing method of LED with a phosphor layer which does not easily come off.\n2. Description of Related Art\nLight-emitting diodes (LEDs) are popular recently. LEDs have advantages such as energy saving, electricity saving, high efficiency, short response time, extensive life span, mercury-free, and environmental protection benefits. Therefore, LEDs are considered as the best light source of new generation illumination. At the present, the phosphor glue of a LED is formed by the processes of forming a reflective cup on the substrate and deposing a LED chip in the reflective cup, and then the reflective cup is filled with phosphor glue via dispensing. However, the connections between phosphor layer and substrate and reflective cup are unstable, whereby the phosphor layer has the problem of easily coming off from the substrate and the reflective cup.\nIn view of the above-mentioned problem, it is necessary to provide a light-emitting diode manufacturing method which the phosphor layer of the light-emitting diode does not easily come off."}
{"text": "Traffic accidents worldwide cause over one million deaths per year, and over 30,000 deaths per year in the U.S. alone. Despite steadily increasing safety standards for automobiles and road construction, distracted driving, intoxicated driving, driver incompetence or inability, dangerous roads and weather conditions, high traffic roads, and extensive road commutes remain as perpetual factors in the lack of a perspicuous decline in traffic-related deaths and injuries. The advent of autonomous vehicle technology—along with the persistent increases in machine learning and artificial intelligence technology—may prove to circumvent many of the unfortunate factors that lead to traffic accidents.\nWidespread concerns regarding autonomous vehicles on public roads typically relate to the ability of such vehicles to make safe and trustworthy decisions when confronted with complex situations. On a typical journey, an autonomous vehicle may encounter countless decision-making instances where loss of life is possible—however unlikely. Road intersections include traffic signaling systems that can range from simple three-bulb faces to complex directional and yielding signals. Safe, reliable, skillful, and responsible decision-making by autonomous vehicles at any intersection is necessary in order to advance the public use of autonomous vehicles and eventually prevent the vast majority of traffic accident types occurring in present road environments."}
{"text": "1. Field of the Disclosure\nEmbodiments described herein relate generally to amine solvent solutions that absorb acid gases and more particularly to additives that decrease the presence of amine-derived contaminants in- and/or degradation of such amine solvent solutions.\n2. Description of the Related Art\nPlants such as refineries, processing plants, industrial plants and the like, may include an amine treating system to treat liquid and/or gas feed streams. Generally, such feed stream treatment includes an amine solvent solution to absorb acid gases from the feed stream. Acid gases include gases such as hydrogen sulfide (H2S), carbon disulfide (CS2) carbonyl sulfide (COS), and carbon dioxide (CO2). Acid gases may later be removed from the amine solvent solution to regenerated and recycle the amine solvent solution for additional use.\nAmine-derived contaminants, however, can accumulate in the amine solvent solution. If left unchecked, these contaminants can have an adverse effect on the amine treating system. For instance, amine-derived contaminants are associated with a decrease in the amine solvent solution's ability to absorb acid gases and an increase in corrosion within the amine treating system.\nGenerally, amine-derived contaminants result from a reaction or association between the amine in the amine solvent solution with another molecule resulting in another contaminant or a reaction intermediate involving a contaminant. These other contaminants/intermediates include acid gases, oxygen, strong anions, carboxylic acids, and others. Contaminants such as acid gases may come from the feed stream being treated, but contaminants may come from any source such as the make-up water for the amine solvent solution or any other source.\nOne type of amine-derived contaminant is heat-stable salts. Heat-stable salts form when a strong anion, such as chloride, formate, or acetate, reacts with or binds an amine cation. The resultant salts are heat-stable because the addition of heat does not readily regenerate the amine solvent solution.\nAnother type of amine-derived contaminant is amine-derived degradation products. Generally, amine-derived degradation products result from the breakdown of amine molecules into a different chemical species. The chemistry of degradation product formation is complex, and in many cases, the reactions are irreversible. A simplified example included the reaction of oxygen or an acid gas with the amine eventually to form an amine-derived degradation product. Alternatively or additionally, oxygen or an acid gas may react with another contaminant to form an intermediate that reacts with the amine to form the amine-derived degradation product. Of course, formation of amine-derived degradation products is not limited to the forgoing, much simplified, examples.\nSince there are many ways in which heat-stable salts and amine-derived degradation products can be produced, they can, and usually are, both be present in an amine treating system at the same time. Furthermore, amine treating systems can tolerate only so much accumulation of such amine-derived contaminants before it must be addressed. There are many different ways to clean an amine treating system once the contaminants are produced, but there remains a need for ways to avoid or to decrease amine-derived contaminants from forming in the first place."}
{"text": "The present invention relates to printing apparatus, and more particularly to an improved apparatus and method which facilitates production of quality images upon repetitive printing operations.\nThe printing of multicolored and other complex images on articles (such as sheets of textiles, plastic or the like) is commonly accomplished by screen printing machines. Generally, these machines are provided with a conveyorized printing blanket driven over a pair of spaced apart rotating drums. The article is placed flat on the blanket and indexed to each in a series of printing stations. At each station, a printing head is lowered onto the article and a printing operation is performed. For example, a print squeegee is stroked across the surface of a horizontal screen in registry with the article so as to force printing ink through the screen and onto the article, thereby effecting printing. In this manner, each in a series of printing steps are performed on the article to obtain a desired complex image.\nTo assure quality, it has been found desirable that the printing head be in precise registration with the article. However, manufacturing imperfections inherent in printing blankets, i.e., the nonuniform center or neutral plane within, have been found to inhibit precise registration. In particular, upon each indexing movement of the blanket (and associated angular displacement of the drums), the distance traveled by the blanket circumference (over the drums) varied with the neutral plane of the blanket. While these variations were small initially, they compounded upon repeated printing operations, visibly offsetting the blanket from a desired position relative to the printing head. This resulted not only in overlapping images and other printing inaccuracies, but also complicated set up by requiring numerous adjustments to obtain acceptable, though seldom precise, registration."}
{"text": "1 . Field\nThe following description relates to a rectifier which may be used with a wireless power receiver.\n2 . Description of Related Art\nResonance power may include electromagnetic energy. A conventional resonance power transferring system may transmit power wirelessly, and may include a source device that transmits a resonance power and a target device that transmits a resonance power. Resonance power may be transferred from the source device to the target device.\nWhen an amount of current increases due to properties of a diode included in a conventional rectifier in a wireless power receiver (i.e., the target device of the wireless power transmission system), a voltage drop may increase due to resistance of the diode.\nVarious products, such as, for example, high-power applications that consume more than 100 W power and low-power applications that consume less than 10 W, have been studied. However, it has been found that for a wireless power transmission system that consumes about 10 W, the total efficiency is low, for instance, only about 60%."}
{"text": "1. Field of the Invention\nThe present invention relates to a recording head cleaning apparatus, an image recording apparatus and a recording head cleaning method, and more particularly to a recording head maintenance system in an image recording apparatus which employs a drum conveyance system to convey a recording medium.\n2. Description of the Related Art\nAs a general image recording apparatus, it is suitable to use an inkjet recording apparatus, which forms a desired image on a recording medium by ejecting and depositing colored inks from a plurality of nozzles provided in an inkjet head. If the inkjet head is operated for a long period of time, adhering matter such as solidified ink or paper dust from the recording medium, and the like, adhere to the nozzle surface. In particular, if adhering matter becomes attached to the vicinity of the nozzles and the nozzle apertures, this gives rise to deflection of the ejection direction of the ink ejected from the nozzles, or reduction in the ejection volume, and so on, and therefore an inkjet recording apparatus is composed in such a manner that cleaning of the nozzle surface is carried out appropriately.\nJapanese Patent Application Publication No. 2000-094703 discloses a cleaning apparatus which applies a cleaning liquid in a non-contact fashion to an inkjet head which is horizontally installed, by rotating an application roller having a cylindrical shape which is immersed in the cleaning liquid. However, the nozzle surface of an inkjet head that records an image on a recording sheet held on the outer circumferential surface of a cylindrical conveyance roller has a prescribed inclination with respect to the horizontal, in order to maintain a uniform distance with respect to the recording sheet. When carrying out cleaning of the inkjet head thus disposed in the inclined state, using an application roller having cleaning liquid held on the surface thereof, it is necessary to incline the application roller in such a manner that the application roller is parallel to the nozzle surface.\nIf an application roller 200 in a cleaning apparatus 202 shown in FIG. 8 is inclined in accordance with the inclination of a nozzle surface 204A of an inkjet head 204, then the state depicted in FIG. 9 is obtained. In the cleaning apparatus 202 depicted in FIG. 9, a liquid surface 212 of a cleaning liquid 210 accommodated in a case 208 is inclined with respect to a rotational axle 216 of the cleaning roller 200, and therefore it is almost impossible to create a coating layer (liquid pool) 214 of the cleaning liquid in an upper portion 200B of an inclined surface 200A of the application roller 200, and hence the coating layer 214 of the cleaning liquid assumes a non-uniform shape and collects in a lower portion 200C of the inclined surface 200A. In the coating layer 214 having an instable shape of this kind, it is difficult to achieve stable application of the cleaning liquid to the nozzle surface 204A of the inkjet head 204."}
{"text": "This section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present disclosure, which are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present disclosure. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nA vehicle that uses one or more battery systems for providing all or a portion of the motive power for the vehicle can be referred to as an xEV, where the term “xEV” is defined herein to include all of the following vehicles, or any variations or combinations thereof, that use electric power for all or a portion of their vehicular motive force. As will be appreciated by those skilled in the art, hybrid electric vehicles (HEVs) combine an internal combustion engine propulsion system and a battery-powered electric propulsion system, such as 48 volt or 130 volt systems. The term HEV may include any variation of a hybrid electric vehicle. For example, full hybrid systems (FHEVs) may provide motive and other electrical power to the vehicle using one or more electric motors, using only an internal combustion engine, or using both. In contrast, mild hybrid systems (MHEVs) disable the internal combustion engine when the vehicle is idling and utilize a battery system to continue powering the air conditioning unit, radio, or other electronics, as well as to restart the engine when propulsion is desired. The mild hybrid system may also apply some level of power assist, during acceleration for example, to supplement the internal combustion engine. Mild hybrids are typically 96V to 130V and recover braking energy through a belt or crank integrated starter generator. Further, a micro-hybrid electric vehicle (mHEV) also uses a “Stop-Start” system similar to the mild hybrids, but the micro-hybrid systems of a mHEV may or may not supply power assist to the internal combustion engine and operates at a voltage below 60V. For the purposes of the present discussion, it should be noted that mHEVs typically do not use electric power provided directly to the crankshaft or transmission for any portion of the motive force of the vehicle, but an mHEV may still be considered as an xEV since it does use electric power to supplement a vehicle's power needs when the vehicle is idling with internal combustion engine disabled and recovers braking energy through an integrated starter generator. In addition, a plug-in electric vehicle (PEV) is any vehicle that can be charged from an external source of electricity, such as wall sockets, and the energy stored in the rechargeable battery packs drives or contributes to drive the wheels. PEVs are a subcategory of electric vehicles that include all-electric or battery electric vehicles (BEVs), plug-in hybrid electric vehicles (PHEVs), and electric vehicle conversions of hybrid electric vehicles and conventional internal combustion engine vehicles.\nMicro Hybrid technology can use a dual voltage architecture, such as a traditional 12V vehicular electrical system used in conjunction with a lead-acid battery, and a 48 volt vehicular electrical system used in conjunction with a Lithium-ion battery. 12 volt electrical system, as used herein, refers to a traditional vehicular electrical system that operates at a nominal 12 volts. The actual voltage varies dynamically depending in part on the charge state of the battery and the load, and an any point in time can be more or less than 12 volts. 48 volt electrical system, as used herein, refers to a vehicular electrical system that operates at a nominal 48 volts, such as one using an LI-ion battery. The actual voltage varies dynamically depending in part on the charge state of the battery and the load, and an any point in time can be more or less than 48 volts. The 12 volt system can include things such as lights, audio/entertainment, electronic modules and ignition. The 48 volts system can include the A/C compressor, active chassis, and regeneration. These systems support higher power loads and provide redundancy. Typically an 8-10 kW motor/generator captures energy for regeneration, supports re-start and supports higher power loads. A DC/DC converter bridges between the higher 48 volt system and the traditional 12 volt system.\nSuch a micro hybrid vehicle can change electrical load management due to high power regeneration, and provide for electrification of new loads such as air conditioning, active chassis and safety, electric supercharging, as well as result in increased fuel efficiency.\nThe DC-DC converter needed for to bridge the systems should be able to provide sufficient power without taking excess space. Moreover, it should be able to withstand the vehicular environment, including high temperatures."}
{"text": "1. Field of the Invention\nThe invention relates to a method of fabricating a dynamic random access memory (DRAM), and more particularly to a method of fabricating a capacitor in a DRAM.\n2. Description of the Related Art\nModern semiconductor fabrication technique in an ultra large scale integration (ULSI) highly increases the circuit density on a chip. The increase of circuit density causes the downsizing of devices and the increase of device packing density. Recently, enhanced resolution of photolithography technique, the development of anisotropic plasma etching and other improvements of semiconductor fabrication have all been advantageous to device downsizing. However, in order to develop towards a further higher circuit density, some breakthrough is required for semiconductor fabrication.\nDRAM is a device broadly used in electronic industry for data storage due to the characteristic of increased circuit density in an integrated circuit (IC). The stored information or message is determined by the charges stored in an internal capacitor of a memory cell. The access of data is performed by operating the read/write circuit and the peripheral memory in a chip. A single DRAM memory cell comprises a field effect transistor (FET) and a capacitor as a bit for representing a binary data.\nAs the number of transistors in a DRAM greatly increases, the dimension of the transistors reduces. Thus, during storing charges, a acceptable signal-to-noise (S/N) ratio is difficult to maintain. By decreasing the charges in a capacitor to enhance the S/N ratio, the refresh cycles for storing charges is correspondingly increased.\nBeing restricted by the limited available surface area of a capacitor in a memory cell, to supply sufficient capacitance to the chip without increasing the occupied space on the substrate, a special and effective capacitor structure is needed to meet the requirement of semiconductor fabrication. As example, a trench capacitor, a cylinder capacitor, and a stack capacitor have been developed and used. However, due to the high complexity of fabrication, the trench capacitor is not as common as the cylinder capacitor and the stack capacitor. The disadvantages of these structures are the complex process and the high cost of fabrication.\nIn FIG. 3a to FIG. 3g, a conventional method of fabricating a cylinder capacitor in a DRAM is shown.\nReferring to FIG. 3a, on a silicon substrate 300 having a metal-oxide-semiconductor (MOS) formed thereon, an oxide layer 301 and a silicon nitride layer 302 are formed in sequence. The silicon nitride layer 302 is used as an etch stop in the subsequent etching process.\nIn FIG. 3b, using photolithography and etching, the silicon nitride layer 302 and the oxide layer 301 are patterned to form an opening 309, so that the silicon substrate 300 is exposed within the opening 309, for example, a doped region in the MOS is exposed. A poly-silicon layer 303 is formed on silicon nitride layer 302 and fills the opening 309.\nIn FIG. 3c, the poly-silicon layer 303 is etched back until the surface of the poly-silicon layer and the surface of the silicon nitride layer 302 are at a same level.\nIn FIG. 3d, an oxide layer 304 is formed over the substrate 300. Using photolithography and etching, the oxide layer 304 is patterned to form an opening 310, so that the poly-silicon layer 303 within the opening 310 and a part of the surface of the silicon nitride layer 302 are exposed. A poly-silicon layer 305 is formed to cover the opening 310 and the oxide layer 304, and thus, the poly-silicon layer 303 and the poly-silicon layer 305 are electrically connected. An oxide layer 306 is formed on the poly-silicon 305.\nIn FIG. 3e, the oxide layer 306 is etched back with the poly-silicon layer 305 as an etch stop. The poly-silicon layer 305 is etched back with the oxide layer 304 as an etch stop.\nIn FIG. 3f, the remaining oxide layer 306 and the remaining oxide layer 304 are removed by wet etching with the silicon nitride layer 302 as an etch stop.\nIn FIG. 3g, an insulation layer 307, for example, an oxide/nitride/oxide (ONO) layer, is formed over the substrate 300. A poly-silicon layer 308 is formed on the insulation layer 307. The fabrication of a conventional cylinder capacitor in a DRAM is formed.\nIn the above method, two photolithography and etching processes are used, so that two photo-masks are required. Thus, the possibility of misalignment is increased, the process is more complex, and the cost of fabrication is high."}
{"text": "1. Field of the Invention\nThe present invention relates generally to cellulose esters and/or ionic liquids. One aspect of the invention concerns processes for producing cellulose esters in ionic liquids.\n2. Description of the Related Art\nCellulose is a β-1,4-linked polymer of anhydroglucose. Cellulose is typically a high molecular weight, polydisperse polymer that is insoluble in water and virtually all common organic solvents. The use of unmodified cellulose from wood or cotton products, such as in the housing or fabric industries, is well known. Unmodified cellulose is also utilized in a variety of other applications usually as a film (e.g., cellophane), as a fiber (e.g., viscose rayon), or as a powder (e.g., microcrystalline cellulose) used in pharmaceutical applications. Modified cellulose, including cellulose esters, are also utilized in a wide variety of commercial applications. Cellulose esters can generally be prepared by first converting cellulose to a cellulose triester, then hydrolyzing the cellulose triester in an acidic aqueous media to the desired degree of substitution (“DS”), which is the average number of ester substituents per anhydroglucose monomer. Hydrolysis of cellulose triesters containing a single type of acyl substituent under these conditions can yield a random copolymer that can comprise up to 8 different monomers depending upon the final DS.\nIonic liquids (“ILs”) are liquids containing substantially only anions and cations. Room temperature ionic liquids (“RTILs”) are ionic liquids that are in liquid form at standard temperature and pressure. The cations associated with ILs are structurally diverse, but generally contain one or more nitrogens that are part of a ring structure and can be converted to a quaternary ammonium. Examples of these cations include pyridinum, pyridazinium, pyrimidinium, pyrazinium, imidazolium, pyrazolium, oxazolium, triazolium, thiazolium, piperidinium, pyrrolidinium, quinolinium, and isoquinolinium. The anions associated with ILs can also be structurally diverse and can have a significant impact on the solubility of the ILs in different media. For example, ILs containing hydrophobic anions such as hexafluorophosphates or triflimides have very low solubilities in water, while ILs containing hydrophilic anions such chloride or acetate are completely miscible in water.\nThe names of ionic liquids can generally be abbreviated according to the following convention. Alkyl cations are often named by the first letters of the alkyl substituents and the cation, which are given within a set of brackets, followed by an abbreviation for the anion. Although not expressively written, it should be understood that the cation has a positive charge and the anion has a negative charge. For example, [BMlm]OAc indicates 1-butyl-3-methylimidazolium acetate, [AMlm]Cl indicates 1-allyl-3-methylimidazolium chloride, and [EMlm]OF indicates 1-ethyl-3-methylimidazolium formate.\nIonic liquids can be costly; thus, use of ionic liquids as solvents in many processes may not be feasible. Despite this, methods and apparatus for reforming and/or recycling ionic liquids have heretofore been insufficient. Furthermore, many processes for producing ionic liquids involve the use of halide and/or sulfur intermediates, or the use of metal oxide catalysts. Such processes can produce ionic liquids having high levels of residual metals, sulfur, and/or halides."}
{"text": "1. Field of the Invention\nThe field of the invention is data processing, or, more specifically, methods, systems, and products for determining a bisection bandwidth for a multi-node data communications network.\n2. Description of Related Art\nThe development of the EDVAC computer system of 1948 is often cited as the beginning of the computer era. Since that time, computer systems have evolved into extremely complicated devices. Today's computers are much more sophisticated than early systems such as the EDVAC. Computer systems typically include a combination of hardware and software components, application programs, operating systems, processors, buses, memory, input/output devices, and so on. As advances in semiconductor processing and computer architecture push the performance of the computer higher and higher, more sophisticated computer software has evolved to take advantage of the higher performance of the hardware, resulting in computer systems today that are much more powerful than just a few years ago.\nParallel computing is an area of computer technology that has experienced advances. Parallel computing is the simultaneous execution of the same task (split up and specially adapted) on multiple processors in order to obtain results faster. Parallel computing is based on the fact that the process of solving a problem usually can be divided into smaller tasks, which may be carried out simultaneously with some coordination.\nParallel computers execute parallel algorithms. A parallel algorithm can be split up to be executed a piece at a time on many different processing devices, and then put back together again at the end to get a data processing result. Some algorithms are easy to divide up into pieces. Splitting up the job of checking all of the numbers from one to a hundred thousand to see which are primes could be done, for example, by assigning a subset of the numbers to each available processor, and then putting the list of positive results back together. In this specification, the multiple processing devices that execute the individual pieces of a parallel program are referred to as ‘compute nodes.’ A parallel computer is composed of compute nodes and other processing nodes as well, including, for example, input/output (‘I/O’) nodes, and service nodes.\nParallel algorithms are valuable because it is faster to perform some kinds of large computing tasks via a parallel algorithm than it is via a serial (non-parallel) algorithm, because of the way modern processors work. It is far more difficult to construct a computer with a single fast processor than one with many slow processors with the same throughput. There are also certain theoretical limits to the potential speed of serial processors. On the other hand, every parallel algorithm has a serial part and so parallel algorithms have a saturation point. After that point adding more processors does not yield any more throughput but only increases the overhead and cost.\nParallel algorithms are designed also to optimize one more resource the data communications requirements among the nodes of a parallel computer. There are two ways parallel processors communicate, shared memory or message passing. Shared memory processing needs additional locking for the data and imposes the overhead of additional processor and bus cycles and also serializes some portion of the algorithm.\nMessage passing processing uses high-speed data communications networks and message buffers, but this communication adds transfer overhead on the data communications networks as well as additional memory needed for message buffers and latency in the data communications among nodes. Designs of parallel computers use specially designed data communications links so that the communication overhead will be small but it is the parallel algorithm that decides the volume of the traffic.\nMany data communications network topologies are used for message passing among nodes in parallel computers. Such network topologies may include for example, a tree, a rectangular mesh, and a torus. In a tree network, the nodes typically are connected into a binary tree: each node typically has a parent and two children (although some nodes may only have zero children or one child, depending on the hardware configuration). A tree network typically supports communications where data from one compute node migrates through tiers of the tree network to a root compute node or where data is multicast from the root to all of the other compute nodes in the tree network. In such a manner, the tree network lends itself to collective operations such as, for example, reduction operations or broadcast operations. The tree network, however, does not lend itself to and is typically inefficient for point-to-point operations. A rectangular mesh topology connects compute nodes in a three-dimensional mesh, and every node is connected with up to six neighbors through this mesh network. Each compute node in the mesh is addressed by its x, y, and z coordinate. A torus network connects the nodes in a manner similar to the three-dimensional mesh topology, but adds wrap-around links in each dimension such that every node is connected to its six neighbors through this torus network. A mesh or a torus network generally lends itself well for point-to-point communications. In computers that use a torus and a tree network, the two networks typically are implemented independently of one another, with separate routing circuits, separate physical links, and separate message buffers. Other network topology often used to connect nodes of a network includes a star, a ring, or a hypercube.\nThe different network topologies mentioned above each have different characteristics that impact the performance for data communications among the nodes in a network. One important performance characteristic is bisection bandwidth. Bisection bandwidth is the total network bandwidth that can be achieved when nodes from different, approximately equal size network partitions communication with one another. Bisection bandwidth may therefore be used as a measure of the bandwidth through a network bottleneck and may be used to represent the effectiveness of a network at handling injected traffic. Networks having a higher value for bisection bandwidth generally handle injected traffic better than networks having a lower value for bisection bandwidth. Because bisection bandwidth is such an effective measurement of how well any given network handles injected traffic, readers will appreciate any advancement in determining a bisection bandwidth for a multi-node data communications network."}
{"text": "1. Field of the Invention\nThe present invention relates to a control for controlling the force applied to splitter jaw clutches during engagement and disengagement under various vehicle operating conditions.\nThe adaptive control may be utilized with a fully or partially automated vehicular transmission system or with the range or, preferably, the splitter section of a controller-assisted, manually shifted transmission.\n2. Description of the Prior Art\nCompound manually shifted mechanical transmissions of the range, splitter and/or combined range/splitter type are in wide use in heavy-duty vehicles and are well known in the prior art, as may be seen by reference to U.S. Pat. Nos. 4,754,665; 5,272,929; 5,370,013 and 5,390,561, 5,546,823; 5,609,062 and 5,642,643, the disclosures of which are incorporated herein by reference. Typically, such transmissions include a main section shifted directly or remotely by a manual shift lever and one or more auxiliary sections connected in series therewith. The auxiliary sections most often were shifted by a slave actuator, usually pneumatically, hydraulically, mechanically and/or electrically operated, in response to manual operation of one or more master switches. Shift controls for such systems by be seen by reference to U.S. Pat. Nos. 4,455,883; 4,550,627; 4,899,607; 4,920,815; 4,974,468; 5,000,060; 5,272,931; 5,281,902; 5,222,404 and 5,350,561, the disclosures of which are incorporated herein by reference.\nFully or partially automated transmission systems wherein a microprocessor-based electronic control unit (ECU) receives input signals indicative of various system operating conditions and processes same according to logic rules to issue command output signals to one or more system actuators are known in the prior art, as may be seen by reference to U.S. Pat. Nos. 4,361,060; 4,593,580; 4,595,986; 4,850,236; 5,435,212; 5,582,069; 5,582,558; 5,620,392; 5,651,292 and 5,679,096; 5,682,790; the disclosures of which are incorporated herein by reference.\nSystems wherein variable force is used to engage auxiliary sections, usually to protect range or splitter section synchronizers, are known in the prior art, as may be seen by reference to U.S. Pat. Nos. 5,186,066; 5,193,410; 5,199,314 and 5,224,392, the disclosures of which are incorporated herein by reference."}
{"text": "Known examples of the antibiotics of anthracycline type include daunomycin (daunorubicin; U.S. Pat. No. 3,616,242) and adriamycin (doxorubicin; U.S. Pat. No. 3,590,028), which both may be obtained from the culture broth of a microorganism of actinomycetes. These two compounds exhibit a wide range of antitumor spectra against a variety of experimental tumors and have been used as a chemotherapeutic agent in the clinic practice. Daunomycin and adriamycin are the compound of the general formula ##STR1## wherein R is a hydrogen atom or a hydroxyl group. Daunomycin (the compound of the formula (a) above where R is the hydrogen atom) and adriamycin (the compound of the formula (a) where R is the hydroxyl group) can show a fairly high antitumor activity against various kinds of tumors. However, these two compounds are not necessarily an antitumor agent which are completely satisfactory. Thus, daunomycin and adriamycin have been shown to exhibit wide antitumor spectra against the experimental tumors and also have widely been used as a valuable agent for the therapeutic treatment of tumors in the clinic practice. On the other hand, it is known that daunomycin and adriamycin can bring about heavy adverse side-effects that they can cause cardiac toxicity and a decrease in the number of leukocytes and falling-off of the hair in the patients who received the administration of these agents. It is reported that the glycoside linkage between the daunosaminyl group of the formula ##STR2## and the hydroxyl group at the 7-position of daunomycinone or adriamycinone is likely to be broken in vivo by the hydrolysis, and that the moiety of the aglycon as formed by the in vivo hydrolysis, namely the daunomycinone or adriamycinone shows a higher cardiac toxicity than the daunomycin or adriamycin itself.\nIn the past, therefore, some researches were already made in an attempt to provide new daunomycinrelated compounds which possess a higher anticancer activity and a lower toxicity than daunomycin and adriamycin. For instance, studies for discovering and producing new daumonycin-analogous compounds by fermentative methods, semi-synthetic methods, total synthetic methods or enzymatical conversion methods were conducted. Such particular compounds previously proposed include, for example, aclacinomycins A and B (F. Arcoamone \"Topics in Antibiotic Chemistry\" Vol. 2, pp. 102-279, published from Elis Horwood Limited, U.S.A.; and U.S. Pat. No. 3,988,315), 4'-O-tetrahydropyranyladriamycin (West Germany Patent No. 2,831,579 and Japanese patent publication No. 47194/81) and N-mono-benzyl- or N-di-benzyl-adriamycin (U.S. Pat. No. 4,177,264).\nFurther, the specification of U.S. Pat. No. 4,427,664 of Horton et al describes a chemical structure of compounds represented by the general formula ##STR3## wherein R.sup.1 is a hydrogen atom and R.sup.2 is a methoxy group; or R.sup.1 is a hydroxyl group and R.sup.2 is a methoxy group; or R.sup.1 and R.sup.2 are each a hydrogen atom; or R.sup.1 is a hydrogen and R.sup.2 is a hydroxyl group, and X is an iodine, chlorine, bromine or fluorine atom and Y is a hydroxyl group or acetoxy group, and which compounds are of such structure that an aglycon selected from the group consisting of daunomycinone, desmethoxydaunomycinone, adriamycinone and carminomycinone is linked through the oxygen atom at the 7-position thereof to the 1'-position of a sugar of a 2'-halo-.alpha.-L-hexopyranose of the .alpha.-L-manno type or .alpha.-L-talo type. The method of producing a compound of the above formula (b) which is described in the specification of U.S. Pat. No. 4,427,664 is the method wherein an aglycon such as daunomycinone and a glycal, for example, 3,4-di-O-acetyl-L-rhamnal or 3,4-di-O-acetyl-L-fucal, which corresponds to the sugar to be linked to said aglycon are dissolved together in substantially equimolar proportions in a mixture of aprotic organic solvents consisting of anhydrous acetonitrile and tetrahydrofuran and wherein to the resulting solution is then added an iodination agent such as N-iodosuccinimide together with a solvating agent such as dichloromethane at a low reaction temperature so that the aglycon reacts with the said glycal in such way that the glycal used is linked to the 7-hydroxyl group of the aglycon with accompanying by an alkoxyhalogenation of the glycal. According to this method of Horton et al, it happens as described in the specification of said U.S. Pat. No. 4,427,664, that the halogen atom as possessed by the halogenation agent employed is introduced into the 2'-position of the sugar moiety of the compound of the formula (b) formed as the reaction product, and that for instance, when N-iodosuccinimide is employed as the iodination agent, the iodine atom is introduced into the 2'-position of the sugar moiety of the reaction product as obtained.\nThe U.S. Pat. No. 4,427,664 specification discloses an Experimental Example in which 3,4-di-O-acetyl-L-rhamnal of the formula ##STR4## where Ac denotes an acetyl group here and also hereinafter unless otherwise stated is reacted with daunomycinone and N-iodosuccinimide, with accompanying alkoxyhalogenation of said rhamnal compound, to produce 7-O-(3,4-di-O-acetyl-2,6-dideoxy-2-iodo-.alpha.-L-manno-hexopyranosyl)daun omycinone, as well as another Experimental Example in which 3,4-di-O-acetyl-L-fucal of the formula ##STR5## is reacted with daunomycinone and N-iodosuccinimide with accompanying alkoxyhalogenation of said fucal compound, to produce 7-O-(3,4-di-O-acetyl-2,6-dideoxy-2-iodo-.alpha.-L-talo-hexopyranosyl)dauno mycinone. However, this U.S. Patent specification does not disclose any further experimental Examples.\nIn the U.S. Pat. No. 4,427,664 specification, the formula (b) appearing therein refers to that X may broadly be an iodine, bromine, chlorine or fluorine atom, but there is not shown any Experimental Examples in which such a compound of the formula (b) where X is the bromine, chlorine or fluorine atom was virtually synthetized. If the one skilled in the art wishes to synthetize a compound of the above formula (b) where X is the bromine, chlorine or fluorine atom, it is expected that in accordance with the method of producing the compound of the formula (b) taught in the U.S. Pat. No. 4,427,664 specification, he will repeat the procedures of the two Experimental Examples as given in said U.S. patent specification using N-bromosuccinimide, N-chlorosuccinimide or N-fluorosuccinimide as the halogenation agent in place of the N-iodosuccinimide employed by Horton et al. Among the above-mentioned three compounds which are expectedly employable as the halogenation agent in place of the N-iodosuccinimide, a substance which is to be represented by the formula ##STR6## and which may be termed as N-fluorosuccinimide is unknown by now, because it is not described before in any literatures with respect to any process of preparing it and with respect to the physical and chemical properties of it, as far as we, the present inventors, have searched numerous literatures. Accordingly, it is evidently considered that a compound which may be termed as N-fluorosuccinimide is an actually unknown substance which was never prepared in the past prior to these days.\nBesides, it is to be noted that the fluorine element which may be and have usually been considered to belong to the class or family of halogens is possessing an extraordinarily different and higher electric negativity than the other halogen elements, iodine, bromine and chlorine, and also that, as be well known, the fluorine element shows the chemical behaviors very much different from those of the other halogen elements. Accordingly, even when it is assumed that the N-fluorosuccinimide will have been prepared by any chemical process, it is highly probable that the properties of the chemical linkage between the fluorine atom and the succinimide group existing in the N-fluorosuccinimide as an imaginable compound would be very much different from the properties of the other sorts of a halo group linking to the succinimide group, and that such fluoro group is too strongly or too weakly linking to the succinimide group. For these reasons, it is very much hardly conceivable that the N-fluorosuccinimide can act as a fluorination agent to transfer its fluoro group into a second compound and bring about the fluorination of the latter compound. Accordingly, it is not presumable that the N-fluorosuccinimide, even if prepared, can serve especially as the fluorination agent required in the method of producing the compounds of the formula (b) which is described in the U.S. Pat. No. 4,427,664 specification.\nIn short, the U.S. Pat. No. 4,427,664 specification describes the chemical structure of the compounds of the above formula (b) where X may broadly be a chlorine, bromine, iodine or fluorine atom. Of the compounds which are designated by the above formula (b), the compounds of the formula (b) where X is a chlorine or bromine atom are shown merely with reference to their chemical structure in the U.S. Pat. No. 4,427,664 specification but were actually not synthetized concretely. Nonetheless, it is admittable that N-chlorosuccinimide and N-bromosuccinimide are a chlorination or bromination agent already known and necessary and available as the halogenation agent in the process of producing such compound of the formula (b) where X=chlorine or bromine, according to said U.S. patent of Horton et al, and therefore it is deducible from the descriptions of the U.S. Pat. No. 4,427,664 specification that the production of such compounds of the formula (b) where X=chlorine or bromine and isolation of such compounds as produced are possible theoretically in accordance with the method of Horton et al as disclosed in said U.S. patent specification. In contrast, however, it is evident that the U.S. Pat. No. 4,427,664 specification does not disclose or teach a process for really producing such compound of the formula (b) where X is the fluorine atom, to such extent that the process would be workable by chemical experts in view of the disclosure of the U.S. Pat. No. 4,427,664, firstly, because the N-fluorosuccinimide which is deemed as necessary as the fluorination agent for the production of the compound of the formula (b) where X=fluorine, according to the disclosed method of Horton et al is a substance which is still unknown up to now and is very much suspicious to be able to act successfully as the necessary fluorination agent for the intended purpose, and secondly, because the U.S. Pat. No. 4,427,664 of Horton et al does nowhere teach how to prepare the N-fluorosuccinimide. Hence, it is worthy to say that such compound of the above formula (b) where X is the fluorine as shown in the U.S. Pat. No. 4,427,664 was a merely imaginary one which was thought by referring to its chemical structure on the papers in the specification of said U.S. patent and of which utility for the intended antitumor agent is very much suspicious. Accordingly, we, do not believe that such special compound of the formula (b) where X is the fluorine atom could be prepared by the chemical experts according to the disclosure of the U.S. Pat. No. 4,427,664.\nFurthermore, in the U.S. Pat. No. 4,427,664 specification, there is described by the inventors, Horton et al that the compounds of the formula (b) exhibit the antitumor activities against mouse blood cancer cell, Leukemia P 388. More particularly, this U.S. patent specification describes such antitumor activity of 7-O-3,4-di-O-acetyl-2,6-dideoxy-2-ikodo-.alpha.-L-mannohexopyranosyl) daunomycinone (nominated as Compound NSC 331,962 by Horton et al) as tested against Leukemia P 388 but does not describe any data of the antitumor activity of 7-O-(3,4-di-O-acetyl-2,6-dideoxy-2-iodo-.alpha.-L-talohexopyranosyl)daunom ycinone (nominated as Compound NSC 327,472 by Horton et al). While, according to an article of Horton et al reported in the \"Carbohydrate Research\" Vol. 136, pp. 391-396 (1985), they obtained experimental results to show that said Compound NSC 331,962 exhibited an antitumor activity that the increase (in %) of survival days of the mice treated, as compared to the mice untreated (control) (namely, T/C, %), was 247% at a dosage of 50 mg/kg of the test compound when the mice as inoculated with Leukemia P 338 were treated by administration of the test compound; and that said Compound NSC 331,962 exhibited an antitumor activity that the increase (in %) of survival days of the mice treated, as compared to the mice untreated (T/C, %), was 196% at a dosage of 25 mg/kg of the test compound when the mice as inoculated with Leukemia L-1210 -were treated by administration of the test compound (a single dose per day, intraperitoneally given for 9 days). Also, the above-mentioned Compound NSC 327,472 experimentally showed such antitumor activities that the increase (%) of survival days of the Leukemia P 388-inoculated mice treated was 172% at a dosage of 12.5 mg/kg to 25 mg/kg of Compound NSC 327,472, whereas the increase (%) of survival days of the Leukemia P 388-inoculated mice treated decreased to 162% at a further increased dosage of 150 mg/kg of the tested compound.\nApart from the above-mentioned researches of Horton et al, we have made studies in an attempt to produce new daunomycin derivatives or adriamycin derivatives which have better antitumor activity and lower toxicity than daunomycin and adriamycin. As a result, we already succeeded to synthetize a few examples of such daunomycin derivative and adriamycin derivative in which the sugar moiety of daunomycin or adriamycin has chemically been modified and which are useful as the antitumor agent. Thus, we reported 4'-O-tetrahydropyranyl-daunomycins and -adriamycins (Japanese patent publication No. 47194/81); and 3'-deamino-3'-morpholino-daunomycins and -adriamycins (Japanese patent application first publication \"KOKAI\" No. 163393/82).\nFurther, we have made another studies to provide new compounds which are derived by chemical modification with a fluoro group of the 3'-position or 2'-position of kanamycin A and kanamycin B of the aminoglycosidic antibiotics. Thus, we have succeeded to synthetize 3'-deoxy-3'-fluorokanamycin A (Japanese patent application No. 161615/84; U.S. patent application Ser. No. 758,819; European patent application No. 85 4015757.7); 3'-deoxy-3'-fluorokanamycin B (Japanese patent application No. 262700/84); and 2',3'-dideoxy-2'-fluorokanamycin A (Japanese patent application No. 263759/84; U.S. Pat. application Ser. No. 807,485; European patent application No. 85 115901.2).\nIn this way, we already obtained many findings and experiences in the fluorine chemistry of sugars through our studies where the fluoro group is introduced into kanamycins of the glycosidic antibiotics. Based on these findings and experiences, we have now succeeded to synthetize as a new compound a 4-O-benzyl-protected derivative of methyl 2,6-dideoxy-2-fluoro-.alpha.-L-idopyranoside represented by the formula ##STR7## through a multi-stage process with starting from L-fucose of the formula ##STR8## Further, we succeeded to synthetize from the sugar compound of the above formula (g) methyl 2,6-dideoxy-2-fluoro-.alpha.-L-talopyranoside of the formula ##STR9## as a new compound, and further synthetize from the compound of the formula (h) 2,6-dideoxy-2-fluoro-.alpha.,.beta.-L-talopyranose of the formula ##STR10## as a new compound and also a 3,4-di-O-protected-2,6-dideoxy-2-fluoro-.alpha.-L-talopyranosyl halide of the formula ##STR11## wherein A' is a hydroxyl-protecting group, particularly an acyl group, especially a lower alkanoyl group such as acetyl or an aroyl group such as benzoyl and Y is a chlorine bromine or iodine atom, for example, 3,4-di-O-acetyl-2,6-dideoxy-2-fluoro-.alpha.-L-talopyranosyl bromide as a new compound.\nThen, we have now succeeded to produce firstly 7-O-(2,6-dideoxy-2-fluoro-.alpha.-L-talopyranosyl)daunomycinone of the formula ##STR12## as a new compound by reacting the 3,4-di-O-protected-2,6-dideoxy-2-fluoro-.alpha.-L-talpyranosyl halide of the above formula (j) with the 7-hydroxyl group of daunomycinone and then removing the residual hydroxyl-protecting groups (A') from the resulting reaction product.\nFurthermore, by converting the 14-methyl group of the compound of the above formula (k) into a hydroxymethyl group (--CH.sub.2 OH) by treatment with a mild oxidizing agent, we have now succeeded to produce firstly 7-O-(2,6-dideoxy-2-fluoro-.alpha.-L-talopyranosyl)adriamycinone of the formula ##STR13## as a new compound. We also have found that the new compound of the formula (k) and the new compound of the formula (l) have excellent antitumor activities and low toxicities and that the glycoside linkage at the 7-hydroxyl group of these new compounds shows a high stability against hydrolysis by acid. Accordingly, the new compound of the formula (k) and the compound of the formula (l) are interesting for use as antitumor agent owing to their low toxicities coupled with their excellent antitumor activities as demonstrated hereinafter. These new compounds of the formulae (k) and (l) have also high antibacterial activities and are useful as antibacterial agent.\nAccordingly, there is provided an anthracycline derivative represented by a general formula ##STR14## wherein R is a hydrogen atom or a hydroxyl group (Japanese patent application No. 282798/85; U.S. patent application Ser. No. 942,773; European patent application No. 86 117 662.6).\nOf the compounds of the general formula (m), 7-O-(2,6-dideoxy-2-fluoro-.alpha.-L-talopyranosyl)daunomycinone of the formula (k) is in the form of red colored solids having a specific optical rotation [.alpha.].sub.D.sup.25 +197.degree. (c 0.02, chloroform-methanol (1:1)). 7-O-2,6-dideoxy-2-fluoro-.alpha.-L-talopyranosyl)adriamycinone of the formula (l) is in the form of red colored solid having a specific optical rotation [.alpha.].sub.D.sup.25 +194.degree.(c 0.01, chloroform-methanol (1:1)).\nWe have confirmed by animal tests that the compounds of formula (m) exhibit significantly high antitumor activities on experimental tumors and that the level of their antitumor activities is much higher than those of daunomycin and adriamycin, coupled with an acceptably low level of toxicities. Some typical tests on experimental animal tumors are given below."}
{"text": "1. Technical Field\nThe disclosed embodiments relate in general to a control method of sound producing, a sound producing apparatus for a portable device, and a portable apparatus.\n2. Description of the Related Art\nA speaker reproducing audible sound is typically limited to its maximum excursion, temperature rising, transducer's non-linearity and its corresponding amplifier induced nonlinearity. In general, a speaker is driven under a rated power for a long period, and under a maximum power in a short period, which are determined by the manufacturer. The driven power restrictions confine the maximum excursion as well as temperature rising of a speaker in a safe range. For a reproduced signal of a large dynamic range without any distortion, the loudness of the speaker is small, especially for a handheld device under the power restrictions. Consequently, a dynamic range compression (DRC) technology has been introduced in a speaker system for decades so as to trade-off the distortion and loudness under the power restrictions."}
{"text": "1. Field of the Invention\nThe present invention relates to a cordless telephone apparatus having an antenna gain control circuit.\n2. Description of the Related Art\nGenerally, a cordless telephone comprises a base station, which is connected to a wire telephone network, and a plurality of movable portable sets. The base station and the portable sets are connected to each other by a wireless or radio network. In the conventional cordless telephone apparatus, in a case where the portable set initiates a call, a predetermined key pad button is pressed in the portable set. A signal for setting a talking channel is then input into a control circuit of the movable portable set. A resultant talking channel setup is transmitted through an antenna via the wireless network to the base station.\nThe base station receives the transmitted radio signals from the portable set, and modulates the received signal. The modulated signal closes a line switch provided in a hybrid circuit via a baseband circuit of the base station. The wire telephone network and the base station are connected to each other by the line switch, and talking can be performed between the base station and the wire telephone over a line wire. The control circuit of the base station transmits a signal to the portable set by wireless, to indicate that the base station and the wire network are connected to each other.\nIn a case where telephone communications using two portable sets A and B are simultaneously performed, if the portable set A is positioned close to the base station and the portable set B is positioned further away from the base station, the base station receives a high field intensity radio signal from the portable set A and a low field intensity radio signal from portable set B. When the receiver of the base station receives the radio signal from the portable set B, its receiving sensitivity for the radio signal from portable set B is relatively the communication quality with the portable set B is accordingly worsened. The signal/noise may be deteriorated to the point that the base station cannot receive radio signals transmitted the portable set B and a connection between the portable set B and the base station will not be made."}
{"text": "Migraine is a chronic condition with recurrent episodic attacks. It is rather unpredictable illness with its characteristics varying among patients. This unpredictability and variability is also observed within migraine attacks observed in a single patient. Among the most distinguishing features of a migraine is a potential disability caused by the accompanying headache and nausea with or without vomiting as well as extreme sensitivity to sound and light (Headache, 39: 720-727 (1999)). Because of the variability and complexity of the condition, effective management of patients suffering from migraines is challenging.\nMigraine headaches which are considered “primary headaches” are about three times more common in women than in men. Geographically, the occurrence of migraine headaches varies significantly and ranges from 1.5% in Southeast Asia to 14% in Western countries (GRIM, Cephalalgia, 12:229 (1992); JAMA, 267:64-69 (1992) and Pharmacoeconomics, 11: 1-10 (Suppl.1)(1997)).\nSystemic administration of anti-migraine and anti-nausea drugs orally to patients has not been very successful, in part because migraine is often accompanied by nausea and the orally administered drugs are vomited before they can take effect. The only viable route of administration for treatment of nausea and/or migraine is the intravenous or another injectable administration. These typically require a visit at the doctor's office or hospital. The failure to successfully treat migraine or nausea is thus based on a delivery method rather than on the drug effectiveness.\nThe vaginal delivery route of drugs through the vaginal mucosa to the uterus and/or to the general circulation has been discovered by inventors and is disclosed, for example, in the U.S. Pat. Nos. 6,086,909, 6,197,327 and 6,572,874 and in a co-pending application Ser. No. 10/600,849 and 10/349,029, all hereby incorporated by reference.\nAs well as the vaginal delivery route described in the above cited patents and application works, there is still some need for improvement, particularly as it concerns an efficacious quantitative drug delivery.\nThe current invention thus concerns an improved transmucosal delivery of anti-migraine and anti-nausea drugs through vaginal mucosa directly to uterus or to the general circulation which is more efficacious due to a more quantifiable sequestration of the drug within the impermeable layer or layers covering a proximal portion of the vaginal device.\nIt is therefore a primary objective of this invention to provide a vaginal device, such as a tampon, tampon-like foam or another vaginal device which is coated, or of which a proximal portion is coated, with a fluid impermeable layer(s) of film, foil, foam or xerogel forming a strip or an attached or removable cap or cup wherein said impermeable layer further comprises a mucoadhesive composition comprising an anti-migraine or anti-nausea drug, or a combination of both, said composition being released and delivered from said layer(s) substantially quantitatively through the vaginal mucosa into the uterus and/or to the general circulation."}
{"text": "The subject matter disclosed herein relates to amusement park attractions, and more specifically, to providing augmented experiences in amusement park attractions.\nAmusement parks or theme parks may include various entertainment attractions in providing enjoyment to guests (e.g., families and/or people of all ages) of the amusement parks. For example, the attractions may include a ride attraction (e.g., closed-loop track, dark ride, thriller ride, or other similar ride), and there may be themed environments along the ride that may be traditionally established using equipment, furniture, building layouts, props, decorations, and so forth. Depending on the complexity of the themed environments, this could prove to be very difficult and time-consuming to setup and replace the themed environment. It may also be very difficult to setup a themed environment that is entertaining for all passengers on the ride. The same themed environment may be appealing to some passengers, but not others.\nIn addition, due to different motion paths and/or different view perspectives of the passengers in the ride vehicle, it may be difficult to provide the same ride experience to all passengers. For example, passengers sitting in the front row may have better view and thus a more immersive ride experience than passengers in the back row. It is now recognized that it is desirable to include attractions where it may be possible to change attraction themes, or to include or remove certain themed features in such attractions in a flexible and efficient manner relative to traditional techniques. It is also now recognized that it may be desirable to provide an immersive and more personalized or customized ride experience for all passengers."}
{"text": "The present invention, in some embodiments thereof, relates to systems and/or methods for image registration and, more particularly, but not exclusively, to systems and/or methods for registration of ultrasound and Computed Tomography CT and/or Magnetic Resonance (MR) images.\nPhysicians may use several imaging modalities when examining the heart of a patient. Each imaging modality may provide better imaging detail of some tissue types and/or tissue functionalities, but poor imaging of other tissue types and/or functionalities. Therefore, several different modalities may be used to gain an overall clinical picture.\nUltrasound may be used to provide dynamic and/or functional information of the beating heart. Ultrasound may also provide images of the function of the cardiac valves. Computed tomography may provide detailed examination of the heart structure, for example, chambers, blood vessels, connective tissue and/or muscle tissues. However, due to limitations in radiation exposure of the patient, fewer samples may be acquired at different stages of the cardiac cycle, yielding high quality images at some time phases and low quality images at other time phases.\nUS and CT imaging modalities may be different in terms of, for example, image artifacts, intensity levels and/or spatial shape differences. The differences may be due to physical processes related to image formation of each modality. Ultrasound images are obtained by echo (i.e., reflection) of acoustic waves. CT images are obtained from X-rays, optionally accompanied by injected contrast agents. The differences in the image modalities raise difficulties in registering the US and CT images.\nVarious attempts have been made to align US and CT images, in order to fuse the information of both modalities and enable better diagnosis, follow up and visualization during treatment procedures. However, the differences in the modalities, for example, image artifacts, intensity levels, and spatial shape differences due to physics related to image formation in each modality, have made image registration challenging.\nHuang et al., “Rapid dynamic image registration of the beating heart for diagnosis and surgical navigation.” IEEE Trans Med Imaging. 2009 November; 28(11):1802-14. discloses “ . . . a rapid two-step method for registering RT3D US to high-quality dynamic 3-D MR/CT images of the beating heart. This technique overcomes some major limitations of image registration (such as the correct registration result not necessarily occurring at the maximum of the mutual information (MI) metric) using the MI metric.”"}
{"text": "The present invention relates to an automatic-reverse tape player, and more particularly to an operation switching mechanism for use in such an automatic-reverse tape player.\nAutomatic-reverse cassette tape players are widely used in automobiles. Many of these automatic-reverse cassette tape players have a common motor for rotating capstans and reel bases for the purpose of saving electric energy. In ordinary cassette tape players, the rotational speed of a reel base varies depending on the diameter of the tape being wound or unwound on the reel base in a play mode in which the tape is fed at a constant speed by a pair of capstans and a pair of pinch rollers. A slip mechanism is associated with the reel base for allowing the tape that is fed at the constant speed by the capstans and the pinch rollers to be wound on the reel base without being slackened or excessively tensioned.\nThe tape transport mechanism of an automatic-reverse cassette tape player has a pair of slip mechanisms associated with the respective reel bases for feeding the tape at a constant speed in opposite directions. When the tape is fed at the constant speed, rotational forces are transmitted from one of a pair of flywheels rotated in different directions by one capstan motor selectively to one of the slip mechanisms associated with the respective reel bases depending on the direction in which the tape is transported. A pair of intermediate rotatable idlers is disposed between the pair of flywheels and the pair of slip mechanisms. The pair of flywheels and the pair of intermediate rotatable idlers make the entire mechanism complex and result in an increased number of parts. Since the intermediate rotatable idlers are spaced from each other, delicate timing has been required to control the switching between the intermediate rotatable idlers in synchronism with the switching between the pinch rollers, and it is tedious and time-consuming to adjust such timing.\nIt is widely known to control the switching between the directions in which the tape is transported, the switching to a fast feed mode, and the loading and ejecting of the tape with a combination of a reversible control motor and cam gears. For example, there is known a switching mechanism for effecting the above switching operations with an intermediate cam gear swingable into mesh with one of the cam gears, each having tooth-free regions, depending on the direction in which the control motor rotates.\nWith the known structure, the cam gears are concentrically stacked one on the other, and hence are of a complex structure as a whole in their transverse direction. Furthermore, the cam gears do not have a large layout freedom as various members of the switching mechanism and the cam gears are coupled to each other in one position. It has been difficult to achieve delicate timing for controlling reciprocating movement of the loading and tape transport mechanisms while the cam gears are reciprocally rotating through about 360.degree.."}
{"text": "The present invention relates to a filter/separator for a vehicle air conditioning system, and more particularly to a filter/separator which functions as a filter and as a separator as well.\nA conventional air conditioning system for vehicles includes a compressor, a condenser with a cooling fan, a filter for filtering water and other undesired articles contained in coolant, an expansion valve, an evaporator, a blower, and a separator installed between the compressor and the condenser for separation of gaseous coolant and liquid coolant to prevent liquid coolant from entering into the compressor. Nevertheless, the lowering of temperature of the coolant leaving the compressor by the cooling fan of the condenser is not efficient. In addition, the filter is installed in an engine hood which is exposed to a high-temperature environment, rendering the temperature in the filter to be higher than atmospheric temperature. Accordingly, the temperature of coolant cannot be efficiently and effectively lowered, resulting in bad air conditioning, inefficient compressor operation, and energywaste.\nThe present invention provides a filter/separator to mitigate and/or obviate the aforementioned problems."}
{"text": "Multicore optical fibers (MCFs) are single fiber strands created with multiple cores to guide light. MCFs are currently being considered for use in telecommunications, data center communications, and high-performance computing applications. It is believed that spatial division multiplexing with MCFs (i.e., using each core as a separate channel) will enable higher fiber bandwidth and channel density than traditional single-core fibers. While fiber manufacturers are presently developing MCF prototypes, numerous commercialization issues remain, including fiber-to-fiber connectorization and fiber-to-photonics interfacing at the ends of each link. Accordingly, there is a need for new optical coupling systems capable of relaying and transmitting optical signals between MCFs and other optical devices for use with such MCFs."}
{"text": "1. Field\nThe aspects of the disclosed embodiment generally relate to substrate processing systems and, more particularly, to accessing internal areas of the substrate processing system.\n2. Brief Description of Related Developments\nGenerally, conventional load locks, for example, in substrate processing systems have a housing that forms an internal chamber. This internal chamber may house substrates being processed within the substrate processing system. Generally, access is provided to this internal chamber through a manually operated or automatically operated atmospheric door and/or a flat removable lid on the top surface of the load lock. Access to the internal chamber through these conventional atmospheric doors or lids is substantially limited and may require sliding or lifting mechanisms to allow for an exchange of the substrates or other suitable payload within the internal chamber.\nIt would be advantageous to provide a load lock or other suitable processing tool with substantially unhindered access for exchanging a payload or substrate located within an internal chamber of the load lock or processing tool without additional mechanical devices such as sliding (e.g. drawers) or lifting mechanisms that move the workpieces substantially in and out of the load lock or other workpiece handling module."}
{"text": "The N-formyl peptide receptor like-1 (FPRL-1) receptor is a G protein-coupled receptor that is expressed on inflammatory cells such as monocytes and neutrophils, as well as T cells and has been shown to play a critical role in leukocyte trafficking during inflammation and human pathology. FPRL-1 is an exceptionally promiscuous receptor that responds to a large array of exogenous and endogenous ligands, including Serum amyloid A (SAA), chemokine variant sCKβ8-1, the neuroprotective peptide human, anti-inflammatory eicosanoid lipoxin A4 (LXA4) and glucocorticoid-modulated protein annexin A1. FPRL-1 transduces anti-inflammatory effects of LXA4 in many systems, but it also can mediate the pro-inflammatory signaling cascade of peptides such as SAA. The ability of the receptor to mediate two opposite effects is proposed to be a result of different receptor domains used by different agonists (Parmentier, Marc et al. Cytokine & Growth Factor Reviews 17 (2006) 501-519).\nActivation of FPRL-1 by LXA4 or its analogs and by Annexin I protein has been shown to result in anti-inflammatory activity by promoting active resolution of inflammation which involves inhibition of polymorphonuclear neutrophil (PMN) and eosinophil migration and also stimulate monocyte migration enabling clearance of apoptotic cells from the site of inflammation in a nonphlogistic manner. In addition, FPRL-1 has been shown to inhibit natural killer (NK) cell cytotoxicity and promote activation of T cells which further contributes to down regulation of tissue damaging inflammatory signals. FPRL-1/LXA4 interaction has been shown to be beneficial in experimental models of ischemia reperfusion, angiogenesis, dermal inflammation, chemotherapy-induced alopecia, ocular inflammation such as endotoxin-induced uveitis, corneal wound healing, re-epithelialization etc. FPRL-1 thus represents an important novel pro-resolutionary molecular target for the development of new therapeutic agents in diseases with excessive inflammatory responses.\nJP 06172288 discloses the preparation of phenylalanine derivatives of general formula:\nas inhibitors of acyl-coenzyme A:cholesterol acyltransferase derivatives useful for the treatment of arteriosclerosis-related various diseases such as angina pectoris, cardiac infarction, temporary ischemic spasm, peripheral thrombosis or obstruction.\nJournal of Combinatorial Chemistry (2007), 9(3), 370-385 teaches a thymidinyl dipeptide urea library with structural similarity to the nucleoside peptide class of antibiotics:\n\nWO 9965932 discloses tetrapeptides or analogs or peptidomimetics that selectively bind mammalian opioid receptors:\n\nHelvetica Chimica Acta (1998), 81(7), 1254-1263 teaches the synthesis and spectroscopic characterization of 4-chlorophenyl isocyanate (1-chloro-4-isocyanatobenzene) adducts with amino acids as potential dosimeters for the biomonitoring of isocyanate exposure:\n\nEP 457195 discloses the preparation of peptides having endothelin antagonist activity and pharmaceutical compositions comprising them:\n\nYingyong Huaxue (1990), 7(1), 1-9 teaches the structure-activity relations of di- and tripeptide sweeteners and of L-phenyl alanine derivatives:\n\nFR 2533210 discloses L-phenyl alanine derivatives as synthetic sweeteners:\n\nWO2005047899 discloses compounds which selectively activate the FPRL-1 receptor represented by the following scaffolds:\n"}
{"text": "I. Field of the Invention\nThis invention relates generally to control of machines and particularly to eliminating relative misalignment of a plurality of driving means spaced along the length of a moveable member, the member being propelled by the driving means to move linearly transverse to its length and the misalignment being measured parallel to the direction of motion.\nII. Description of Related Art\nAn example of machine construction giving rise to misalignment of driving means of a machine member, referred to as skew, is illustrated in FIG. 1. Machine member 20 is moveable along rails 24 and 26 as indicated by the double ended arrow \"X\". Skew is defined as a difference of the positions of the driving means or driven points of machine member 20 as measured relative to a common reference and parallel to the axis of motion. For example, in FIG. 1 skew is represented by a difference in the values of X' and X\". Movement of member 20 is imparted by, respectively, actuators 28 and 30 driving through transmissions or gear trains 32 and 34, respectively, at distal ends 21 and 22. It will be appreciated that machine constructions may include additional driving means spaced along the length of member 20. The driving means may include, for example, pinions engaging racks on the rails to propel member 20. Other driving devices may be included in the driving means for driving member 20 such as, for example, nuts rotated by actuators and engaging screws mounted parallel to the rails.\nMember 20 and driving means 28 and 30 are arranged to propel member 20 linearly along the rails 24 and 26, the rails being substantially parallel to one another and substantially perpendicular to member 20. Skew may introduce binding of member 20 against rails 24 and 26 as well as poor engagement of the driving means. Skew can result in unsatisfactory response of actuators 24 and 26 to the control thereof and may result in excessive mechanical wear or damage to member 20, rails 24 and 26 or the driving means. Skew may arise as a result of forces acting on member 20 during periods when actuators 28 and 30 are not energized and therefore not active to maintain positions of ends 21 and 22. In light of the adverse consequences of skew, it is desirable to eliminate skew before initiation of extensive motion of member 20.\nIt is known from U.S. Pat. No. 4,045,660 to return a machine member to a desired location following power interruption using values of measured position determined upon power loss and upon power restoration. It is known from U.S. Pat. No. 4,484,287 to restore a moveable machine member to a desired position following power interruption by storing position information in a nonvolatile memory upon power loss for recall upon restoration of power. From U.S. Pat. No. 5,013,988 it is known to use presettable counters in conjunction with absolute encoders to create and update absolute position data, the data being updated so long as power is applied to the measuring system. From U.S. Pat. No. 4,629,955 it is known for control of motion of a machine member driven at distal ends to vary servomechanism gain of the driving means to eliminate skewing therebetween. The control of this patent provides only for reduction of skew attributable to differences in loading on the driving means during execution of motion. Servomechanism position control effected in accordance with this patent does not provide for detection, reduction or elimination of skew between the driving means upon application or restoration of power.\nThe controls of the aforesaid references addressing restoration of position after power interruption all rely on absolute position information which is immediately available upon application or restoration of power. However, in the event sufficient absolute position data is not available on application of power, it must be determined from position transducers. Sufficient position data will not exist if, during a period when position control is disabled, positional changes are not monitored and the range of position measurement transducers is less than the potential magnitude of positional change or the output of the position transducers while the machine member is stationary do not provide any indication of position whatsoever. Position transducers of the latter type include incremental or semi-absolute encoders, which are favored for providing high resolution position measurement over extended distances."}
{"text": "1. Field of the Invention\nThe present invention relates to an exposure apparatus and a method for producing a device in which a substrate is exposed with a pattern via a projection optical system and a liquid.\n2. Description of the Related Art\nSemiconductor devices and liquid crystal display devices are produced by the so-called photolithography technique in which a pattern formed on a mask is transferred onto a photosensitive substrate. The exposure apparatus, which is used in the photolithography step, includes a mask stage for supporting the mask and a substrate stage for supporting the substrate. The pattern on the mask is transferred onto the substrate via a projection optical system while successively moving the mask stage and the substrate stage. In recent years, it is demanded to realize the higher resolution of the projection optical system in order to respond to the further advance of the higher integration of the device pattern. As the exposure wavelength to be used is shorter, the resolution of the projection optical system becomes higher. As the numerical aperture of the projection optical system is larger, the resolution of the projection optical system becomes higher. Therefore, the exposure wavelength, which is used for the exposure apparatus, is shortened year by year, and the numerical aperture of the projection optical system is increased as well. The exposure wavelength, which is dominantly used at present, is 248 nm of the KrF excimer laser. However, the exposure wavelength of 193 nm of the ArF excimer laser, which is shorter than the above, is also practically used in some situations. When the exposure is performed, the depth of focus (DOF) is also important in the same manner as the resolution. The resolution R and the depth of focus δare represented by the following expressions respectively.R=k1·λ/NA  (1)δ=±k2·λ/NA2  (2)\nIn the expressions, λ represents the exposure wavelength, NA represents the numerical aperture of the projection optical system, and k1 and k2 represent the process coefficients. According to the expressions (1) and (2), the following fact is appreciated. That is, when the exposure wavelength λ is shortened and the numerical aperture NA is increased in order to enhance the resolution R, then the depth of focus δ is narrowed.\nIf the depth of focus δ is too narrowed, it is difficult to match the substrate surface with respect to the image plane of the projection optical system. It is feared that the margin is insufficient during the exposure operation. Accordingly, the liquid immersion method has been suggested, which is disclosed, for example, in International Publication No. 99/49504 as a method for substantially shortening the exposure wavelength and widening the depth of focus. In this liquid immersion method, the space between the lower surface of the projection optical system and the substrate surface is filled with a liquid such as water or any organic solvent so that the resolution is improved and the depth of focus is magnified about n times by utilizing the fact that the wavelength of the exposure light beam in the liquid is 1/n as compared with that in the air (n represents the refractive index of the liquid, which is about 1.2 to 1.6 in ordinary cases).\nHowever, the conventional technique as described above involves the following problem. The exposure apparatus, which is disclosed in International Publication No. 99/49504, is constructed such that the liquid is supplied and recovered to form the liquid immersion area on a part of the substrate. In the case of this exposure apparatus, for example, when the substrate stage is moved to the load/unload position in order to unload the substrate having been placed on the substrate stage and load a new substrate in a state in which the liquid in the liquid immersion area is not recovered sufficiently after the completion of the liquid immersion exposure, there is such a possibility that the liquid, which remains on (adheres to) the end portion of the projection optical system, the liquid supply nozzle, and/or the liquid recovery nozzle, may fall onto surrounding units and members including, for example, the guide surface of the stage and the reflecting surface for the interferometer for the stage.\nFurther, when the liquid remains on the optical element disposed at the end portion of the projection optical system, the remaining liquid leaves any adhesion trace (so-called water mark) on the optical element disposed at the end portion of the projection optical system after the evaporation of the remaining liquid. There is such a possibility that any harmful influence may be exerted on the pattern to be formed on the substrate during the exposure process to be subsequently performed. It is also assumed that the liquid immersion area is formed during any process other than the exposure process, i.e., when the reference mark member and/or the reference plane member arranged around the substrate on the substrate stage is used. In such a situation, there is such a possibility that the liquid in the liquid immersion area cannot be recovered sufficiently, the adhesion trace may remain on the member as described above, and the liquid remaining on the member as described above may be scattered."}
{"text": "The present invention provides a process for the rapid and efficient production and examination of polyurethane (“PU”) foam-forming formulations in which only a small amount of material is required.\nScreening of PU foam formulations is generally carried out by hand. Laboratory packets containing from 200 to 300 g of foam are produced after all of the ingredients have been manually weighed, mixed together in a bench stirrer, and the mixture has been poured into paper packets. Disadvantages of this manual screening are the low maximum throughput of 15 packets per day per technician, poor reproducibility resulting from the non-documenting of errors/deviations in weighing, stirring times, stirrer speed, etc., and a laborious determination by hand of reaction parameters such as the cream time, full rise time, fiber time and tack-free time.\nA problem when conducting physical testing of PU foams is that the removal of a plurality of identical sample bodies in accordance with the DIN standard (sample size at least 125 cm3) is virtually impossible due to flow distance phenomena and fluctuations in density (up to 10%) and to the limited sample quantity from one identical batch. Additionally, destructive testing techniques frequently also mean that the sample body can be used only for a single measurement thereby necessitating the use of a plurality of packets which are as nearly identical as possible but which may frequently have properties which differ from one another (for example, differences in densities or in open cell content). Defined storage times must be observed before the samples are examined, in order to avoid or standardize ageing of the samples due to cell gas exchange."}
{"text": "1. Field of Invention\nThe present invention relates to a quick joint for liquid tight joining of lines under high pressure and comprising a sleeve shaped female part and a male part being introducible into said female part, whereby a male part related to the female part comprises a cylindrical part being introducible into the female part the envelope surface of which presents a groove having a certain width, which groove receives a ring being axially displaceable between the two walls of the groove, which ring has substantially the same thickness as the depth of the groove, and which female part comprises a cylindrical space for receiving the cylindrical part of the male part, the envelope surface of said space is provided with at least a groove for receiving elastically movable locking means out off and into the cylindrical part, which means in a locking position cooperates with the edge of the groove being situated closest to the outer end of the male part and which locking means are displaceable out off their locking position by an movement of the ring towards said outer end edge, which movement is made possible by pressing the male part into the female part until the locking means are brought into a frictional engagement with the outer side of the ring. In particular the invention relates to the female part of such a quick joint.\n2. Description of Related Art\nA quick joint according to above is known from EP 0 375 674 and is provided with a ring of spring steel. Such a ring of spring steel provides for a necessary locking at normal hydraulic applications but there is a risk that it will become deformed at extremely high hydraulic pressures. If the locking ring is given larger dimensions in order to withstand higher loads it will become much harder to mount the ring in its groove within the female joint part.\nPCT/SE96/00593 discloses a quick joint of the type given above, as well, wherein a ring shaped flange is coaxially arranged on the outer side of the groove of the envelope surface for receiving locking means being elastically movable out off and into the cylindrical space, and wherein the locking means comprises at least three arc shaped segments arranged in said groove which segments are forced into a an angular position by an elastic element to take a locking position in relation to the symmetrical longitudinal axis of the female part, in which angular position the segments together takes the shape of a truncated cone which supports with its broader base against the outer side of the groove of the female part, radially on the outer side of the flange.\nThese known quick joints become in order to be tight, torsional rigid, which means that those hoses which are connected to the joint will take much of optional rotational movements, such as e.g., at hydraulically operated timber-cutting apparatuses where a front cutting head is rotated in several different planes to facilitate cutting, delimbing and lumber distribution."}
{"text": "1. Field of the Invention\nThe present invention relates to a shift clock generator for generating a shift clock, a timing generator for generating predetermined timing and a test apparatus for testing electronic devices.\n2. Related Art\nConventionally, a test apparatus for testing electronic devices such as a semiconductor device is provided with a timing generator for generating predetermined timing. For example, the test apparatus supplies a test pattern to the electronic device with the timing generated by the timing generator. The timing generator generates the predetermined timing by receiving a reference clock and by delaying the reference clock by a predetermined time.\nThe timing generator has a variable delay circuit section for receiving the reference clock and for delaying the reference clock by the predetermined time and a linearize memory for controlling a value of delay in the variable delay circuit section for example. The variable delay circuit section has a plurality of delay elements in general. The linearize memory stores a delay preset value corresponding to linearization of the predetermined value of delay in the variable delay circuit section. Based on the data stored in the linearize memory, the variable delay circuit section delays the reference clock by passing the reference clock through a route of predetermined delay elements. Although the data stored in the linearize memory is set in advance by design information of the plurality of delay elements, an error occurs between the value of delay in the variable delay circuit section and the delay preset value which is the predetermined value of delay due to dispersion in manufacturing the plurality of delay elements and to ambient temperature in using the delay elements for example.\nConventionally, in order to compensate the error, a shift clock having a phase which is different from that of the the reference clock by a predetermined value is generated and the shift clock is outputted to the outside to measure a compensation value of the value of delay of the shift clock by using a measuring instrument and to linearize the value of delay. The shift clock is compared with the output of the variable delay circuit section to detect the error of the value of delay and to select the data to be stored in the linearize memory based on the error.\nPresently, the present applicant is unaware of related patent documents, so that description thereof will be omitted here.\nConventionally, in order to generate the shift clock having the predetermined phase difference from the reference clock, pulses are inputted to the shift clock to phase-shift the shift clock by a method as described later in connection with FIGS. 3 and 4. Conventionally, the shift clock is phase-shifted by counting pulses of the shift clock and by inserting insertion pulses per predetermined count. However, there is a case when the phase shift amount of the shift clock does not change linearly with respect to the number of insertion pulses and this method causes an error in the phase shift amount of the shift clock when the phase shift amount of the shift clock does not change linearly with respect to the number of insertion pulses.\nThere is also a method of using a memory for storing a number of pulses to be inserted per predetermined phase shift amount in order to eliminate such error. However, in order to accurately measure the value of delay in the variable delay circuit section, resolution of the phase shift amount must be increased and a memory having a wide range of addresses is required. Still more, the number of pulses to be inserted must be stored in each address. Because the number of pulses to be inserted is normally around one to several thousands and such memory must have several tens bits in each address, a memory having a large capacity is required."}
{"text": "1. Field of the Invention\nThe invention relates to a seat carried on a vehicle such as an automobile, an airplane, a ship and an electric train.\n2. Description of Related Art\nAs shown in FIG. 8, in a front seat for an automobile, in order to prevent small items from falling from a gap between a seat cushion 110 and a central control station C towards an action area A of a lift mechanism (not shown) of the seat cushion 110, a cover member 117 hanging towards the lower side of the seat cushion 110 is provided on the side portion of the seat cushion 110 so as to cover the gap (see Japanese Patent Application Publication No. 2013-226932 (JP 2013-226932 A)).\nIn this case, the seat cushion 110 comprises a cushion frame 111 and a cushion cover 115 which covers an outer side of a cushion frame 111 that serves as a base body of the vehicle seat such that a soft cushion (not shown) is interposed between the cushion frame 111 and the cushion cover 115. The cushion cover 115 is integrated with the cushion frame 111 by engaging a hook 116 with a lower end portion 111a of the cushion frame 111. At this time, the lower end portion 111a of the cushion frame 111 is located lower than the position of the cushion cover 115 on the side of the central control station C, and is located on the inner side of the seat cushion 110. Thus, the position of the hook 116 also is the same position, and the lower end portion of the cushion cover 115 is bent to be separated from the central control station C. Thus, the action area A of the lift mechanism on the lower portion of the seat cushion 110 is opened on the side of the central control station C, and small items can fall easily. As a countermeasure, in the invention of JP 2013-226932 A, the action area A is closed on the side of the central control station C by stitching the cover member 117 on the cushion cover 115.\nHowever, according to the vehicle seat described in JP 2013-226932 A, since the cover member 117 is newly provided, there exists a problem that the manufacture cost rises. For example, the time for trimming the cover member 117 is unnecessarily spent. Moreover, the time for stitching the cover member 117 on the cushion cover 115 is also unnecessarily spent."}
{"text": "1. Field of The Invention\nThis invention relates to air conditioner filters.\nMore specifically this invention relates to electrostatic air conditioner filters.\n2. Prior Art\nPrior art shows several filters incorporating layers and electrostatic filtering.\nThis technology dates back for some time. Most of the patents related to electrostatic filtering are improvements on the technology.\nU.S. Pat. No. 4,904,288 to d'Augereau shows electrostatic filtering utilizing a waved steel mesh similar to aluminum screening in order to create turbulence and pockets for dust to collect. d'Augereau and patents cited therein show multiple layer electrostatic and regular air filters.\nThe egg crate design of the polypropylene layers described in this specification are also well known in the art in U.S. Pat. No. 2,724,457.\nOne problem with the prior art is that the utilization of electrostaticly charged elements involves the excessive dependence on layers impeding the air flow. This reduces the effectiveness of the air conditioner, increases energy use and results in difficulty in cleaning the unit.\nAnother problem in the prior art is the use of metallic elements grounding and reducing the effectiveness of the electrostatically charged element.\nIt is therefore one object of the invention to provide for an electrostatic filter made with non-metallic parts so that the electrostatic element retains a greater charge and is not grounded.\nIt is another object of the invention to provide an electrostatic air filter allowing for greater air flow with less resistance and therefore greater efficiency for the air conditioning unit and less strain on the air conditioning motor.\nIt is another object of the invention to provide a layered air filter which is both efficient and easy to clean for reuse.\nThese and other objects of the invention may be more readily observed from the accompanying drawings and detailed description given below."}
{"text": "1. Field of the Invention\nThis invention relates to health related data analysis and more particularly relates to an apparatus system and method for rapid cohort analysis.\n2. Description of the Related Art\nMost corporations, including health insurance corporations, maintain a high volume of data. Such data may be analyzed and exploited for valuable information regarding business trends, and other important statistics. Data mining is a common strategy for identifying and analyzing such data.\nThere are many various forms of data mining. Custom analytic operations may be developed to meet specific needs. Alternatively, commercially available statistical analysis tools, such as Statistical Analysis Software (SAS) may be used to identify statistical trends in data.\nHealth insurance companies typically maintain databases of health insurance claim information, demographic information, and other data about health insurance plan members. Such information may be used to gain valuable insights into disease causes, progressions, and potential cures. Unfortunately, typical methods for analyzing such data are often cumbersome, costly, and require unworkably high processing times and resources.\nThe referenced shortcomings are not intended to be exhaustive, but rather are among many that tend to impair the effectiveness of previously known techniques disease management; however, those mentioned here are sufficient to demonstrate that the methodologies appearing in the art have not been satisfactory and that a significant need exists for the techniques described and claimed in this disclosure."}
{"text": "The present invention is directed to emitting biofluids from drop ejection units, and more particularly to priming mechanisms used to obtain proper drop ejection sensing and controlling the level of biofluid within drop ejection devices.\nVarious designs have been proposed for the ejection of biofluids which permit the high-speed printing of sequences and arrays of drops of biofluids to be used in various tests and experiments. In the present discussion, a biofluid, also called a reagent, may be any substance used in a chemical reaction to detect, measure, examine or produce other substances, or is the substance which is to be detected, measured, or examined.\nBiofluid ejection devices find particular utility in the depositing of drops on to a substrate in the form of a biological assay. For example, in current biological testing for genetic defects and other biochemical aberrations, thousands of the individual biofluids are placed on a glass substrate at different well-defined locations. Thereafter, additional depositing fluids may be deposited on the same locations. This printed biological assay is then scanned with a laser in order to observe changes in the biofluid property.\nIt is critical in these situations that the drop ejection device not be a source of contamination or permit unintended cross-contamination between different biofluids. Also, due to the high cost of these biofluids, and the importance of positioning properly formed drops at highly precise locations, it is important that the drop ejectors operate correctly at the start of the drop ejection process.\nIn view of the foregoing, it has been considered desirable to provide priming mechanisms which ensure the proper delivery of biofluids to an ejector device in a timely, useful manner."}
{"text": "Reinforcements provide structural support without a significant increase in cost and weight. For instance, reinforcements may be used in automobiles to reinforce cavities formed by various parts of the automobile such as a pillar, bumper, etc. To properly transfer loads from one side of the structure to the other, the reinforcement may have features that generally match the inner surfaces of the cavity in which the reinforcement is placed.\nReinforcements may be provided with an adhesive or bonding material that secures the reinforcement within a given cavity. Generally, such materials are provided on outer surfaces of the reinforcement in order to engage corresponding surfaces of the cavity upon insertion of the reinforcement into the cavity. However, such materials may be easily damaged prior to assembly, e.g., during shipping or handling of the reinforcement. Additionally, adhesive materials may be relatively soft, tacky, or otherwise difficult to handle directly, resulting in added difficulty in handling and/or installing the reinforcement."}
{"text": "Numerous papillomavirus sequences were determined, see the publications incorporated herein by reference: HPV-6: de Villiers et al., J. Virology, 40 (1981); HPV-11: Dartmann et al., Virology 151, 124-130 (1986); HPV-16: Seedorf et al., Virology 145, 181-185 (1985); HPV-18: Cole and Danos, Journal of Molecular Biology 93, 599-608 (1987); HPV-31: Goldsborough et al., Virology 171, 306-311 (1989); HPV-33: Cole and Streeck, J. Virology, 58, 991-995 (1986); HPV-54: Favre et al., J. Cancer 45, 40-46 (1990); HPV-56: Lõrincz, J., Gen. Virol. 70, 3099 (1989).\nDetecting and typing of HPV is reported in a number of publications, besides Soutern blotting and other hybridisation techniques, the most widely used techniques are the PCR-based methods, since these methods simultaneously provide high sensitivity, specificity and the flexibility of the assay gives more control to comply the analytical requirements.\nHuman papillomavirus, a member of the Papillomaviridae family, is a DNA tumorvirus, with an 8000 bp of circular genome. The virus shows strong epithelial tropism, and proliferates only in differentiated epithelial cells. The papillomavirus has suspected etiologic role in many different human diseases, for example in different skin diseases, i.e. in verruca, condyloma acuminatum and skin tumours and in other conditions, such as cervical carcinoma, anogenital carcinomas, laryngeal carcinoma. It is well established that the human papillomavirus shows strong correlation with the incidence of these tumors, and this is even true for the pre-cancerous lesions (ClN, VlN, VAIN, PIN, PAIN). HPV can be detected in 99% of the cervical carcinoma patients. This close statistical relationship is possibly caused by the causal role of the HPV in the formation of cervical carcinoma. On the basis of the epidemiological data, the patients to be infected by different HPV genotypes do not have the same level of risk to develop cervical carcinoma. According to these findings the genotypes are classified into low risk, medium risk and high risk classes, and besides these there are not-classified genotypes too. Since the risks are grossly different and the incidence of the HPV infection is very high, the determination of genotype is of great importance.\nThe HPV virus can not be cultivated. The serological diagnosis of HPV infection is limited to detect the exposure to the virus (past or present infection), but can not exactly identify the genotype, the role is mostly limited to epidemiological investigations.\nFor papillomaviruses, exact serologic classification (serotyping) does not exist genotyping is the widely accepted classification method. These can be divided into two groups, according to whether detection is preceded by amplification or not. In one embodiment of the latter method, full length genomic RNA probes are used to detect the denatured HPV DNA genomes, and the heteroduplex is detected with specific antibodies (Hybrid Capture—Digene). According to another method, Southern blot technique is used for detection and genotyping the HPV genotypes. The disadvantage of these methods is the relative insensitivity and partial lack of specificity. In the case of the Hybrid Capture method many publications report different cross-reactions, causing false positive reactions in clinical conditions. The authors reported that the cross-reactions were acceptable only with a cut-off control of high (1 ng/ml) DNA concentration, which underlines the non-desirable coupling between sensitivity and specificity.\nBy the amplification methods this problem does not appear, since the reaction responsible for the sensitivity (amplification) is carried out separately.\nGenerally the amplification techniques differ in the selected amplified genome segment, number of primers, and the applied detection technique. The most frequently used primers are the GP5+-GP6+, MY9-MY11 and the different type-specific PCR reactions.\nThe most frequently used detection techniques are the sequence-specific hybridisation, restriction fragment length polymorphism (RFLP) and the line probe assay (LiPA). Besides these ones, sequencing of amplicons and thymidine pattern generated by dUTP incorporation is used, but less frequently.\nThe analytical characteristics of the amplification techniques vary in wide ranges. The methods can be characterized by the amplifiable genotypes, the analytical sensitivity of the genotype amplification and the specificity and reliability of the detection. In this field the MY9-MY11 degenerated primer system is considered to be the reference reaction. In case of the MY9-MY11 system LiPA hybridisation detection system exists (Innogenetics). The major drawback of the MY9-MY11 system it is difficult to control the degenerated synthesis of the primers that is why the relative ratio of the primer species produced in the synthesis is varying from synthesis to synthesis, which can result in the unpredictable changes of the analytical behaviour of the PCR reaction; secondly, this reaction can amplify the fewest types, compared to the other widespread used reactions. It is well known from the literature, that the system can amplify genotype 51 only in that case, if the HPV genotype 51 type-specific primers are added to the reaction. Using degenerate primer synthesis the relative ratio of the primer species can not be changed, and it is impossible to tailor the primer ratios to achieve better analytical performance and a balanced amplification of genotypes.\nThe GP5+-GP6+ reaction solves the problem only by the use of two carefully selected pair of primers—optimised to the genital HPV sequences—the two primer systems are easy to manage, however the flexibility is lower. The GP5+-GP6+ system can amplify a lot of known HPV genotypes, but the analytical characteristics of the system are not optimal (sensitivity is not balanced with different genotypes), and the two primer approach is constrained in optimisation, e.g. balancing the detection sensitivities for the individual genotypes is highly problematic (except the limited optimisation of the melting temperature and concetration of the MgCl2). It is difficult to adapt the GP5+-GP6+ system to the amplification of other genotypes, which in any case influence its future application, since the need for detecting new genotypes permanently occur. The identification of the genotypes is not solved adequately.\nAnother well known wide genotype-specific amplification method is the L1C method: two-primer system, with two versions, one is using (with the LC1 primer) the L1C2 or the new L1C2 primer, to amplify further genotypes. The detailed description of the L1C amplicon can be found in the literature [Jpn. J. Cancer Research 82, 524-531 (1991)].\nBasically two criteria must be fulfilled by the detection postamplification methods: routine diagnostic applicability (simplicity, costs, time), and the requirement of power of discrimination suitable level of discrimination power. A significant group of methods are not suitable in terms of power of discrimination discrimination power. Therefore the application of the RFLP is limited, because of the short amplified regions, there are not enough diagnostic restriction sites, so often the genotypes can only be classified into groups. Another example the SSCP technique is difficult to refer the complex patterns of the SSCP to genotypes, and also, the robustness of these reactions is not satisfactory, either.\nThe power of discrimination is especially important from the diagnostic point of view to fulfil the requirements of the regulatory authorities. From the aspects of the simplicity and the power of discrimination sequencing is the ideal approach, since its automation is solved and able to detect each genotypes (or even subtypes thereof), if the sample is not a mixture of genotypes. But in the practice it is not widely used, because it is expensive and time-consuming, and its application in routine diagnostic laboratories is not acceptable, and in case of mixed samples none of the genotypes can be determined.\nThe advantage of the hybridisation methods is that their power of discrimination or stringency can easilybe changed, since several parameters of the reaction can be varied in wide ranges, and some forms are easily automated, the reaction is less expensive, and in case of parallel implementation (with some forms) even the time needed is insignificant.\nTherefore there is a need for a new HPV amplification/detection method, which eliminates the disadvantages of the current methods, and it is cheap, easy to reproduce and automate. The invention describes an amplification and hybridisation assay, in which the primers are independently synthesized molecules, therefore their relative ratio can easily be controlled and optimised, and the amplification has a balanced sensitivity. Hybridisation reactions carried out in highly parallel manner comply with the criteria of a low cost, fast, flexible and automatable reaction."}
{"text": "1. Technical Field\nThe present invention relates to position detecting devices and electro-optical devices, and more particularly, to configurations of devices that obtain position information regarding an object to be detected on the basis of values which have been optically detected.\n2. Related Art\nIn typical display devices equipped with an electro-optical device such as a liquid crystal display body, a lighting device such as a backlight may be incorporated in order to make the display screen visible or enhance the visibility. In addition, in the display devices, the display screen may be provided with a pointed position detecting unit such as a touch panel. In this case, when a given position on the display screen is pointed by a pen, a finger or the like, the pointed position is detected and input into an information processing device or the like.\nAs the pointed position detecting unit (positional coordinate inputting unit) such as a touch panel, electrostatic capacitive type touch panels, resistive film type touch panels or the like that mechanically or electrically detect the state of a contact on the display screen are known. In addition, optical touch panels are known in which a grid of infrared light beams is formed across the display screen and photosensors are correspondingly provided so as to detect the infrared light beams, so that, when a finger or the like interrupts infrared light beams, the positional coordinates of the finger or the like can be detected. In general, various types of optical touch panels are known, and examples of the optical touch panels include the ones disclosed in JP-A-2004-295644 and JP-A-2004-303172.\nHowever, in the optical touch panels described above, it is necessary to arrange, adjacent to the display screen, multiple light sources and photosensors, optical switches, or light guide structures or the like in order to support the resolution of positional coordinates to be detected. This increases the number of optical elements, whereby there is a problem in that the manufacturing cost is increased and more power is consumed."}
{"text": "a. Field\nThe instant disclosure relates to steering mechanism or actuators for steerable medical devices. In particular, the instant disclosure relates to actuators in steerable medical devices employing one or more pull wires to deflect a portion of such medical devices in at least one direction and, preferably, in at least two directions.\nb. Background Art\nElectrophysiology catheters are used in a variety of diagnostic, therapeutic, and/or mapping and ablative procedures to diagnose and/or correct conditions such as atrial arrhythmias, including for example, ectopic atrial tachycardia, atrial fibrillation, and atrial flutter. Arrhythmias can create a variety of conditions including irregular heart rates, loss of synchronous atrioventricular contractions, and stasis of blood flow in a chamber of a heart, which can lead to a variety of symptomatic and asymptomatic ailments and even death.\nTypically, a catheter is deployed and manipulated through a patient's vasculature to the intended site, for example, a site within a patient's heart. The catheter typically carries one or more electrodes that can be used for cardiac mapping or diagnosis, ablation, and/or other therapy delivery modes, or both, for example. Once at the intended site, treatment can include, for example, radio frequency (RF) ablation, cryoablation, laser ablation, chemical ablation, high-intensity focused ultrasound-based ablation, microwave ablation, and/or other ablation treatments. The catheter imparts ablative energy to cardiac tissue to create one or more lesions in the cardiac tissue. These lesions disrupt undesirable cardiac activation pathways and thereby limit, corral, or prevent errant conduction signals that can form the basis for arrhythmias.\nTo position a catheter within the body at a desired site, some type of navigation must be used, such as using mechanical steering features incorporated into the catheter (or an introducer sheath). In some examples, medical personnel may manually manipulate and/or operate the catheter using the mechanical steering features.\nIn order to facilitate the advancement of catheters through a patient's vasculature, the simultaneous application of torque at the proximal end of the catheter and the ability to selectively deflect the distal tip of the catheter in a desired direction can permit medical personnel to adjust the direction of advancement of the distal end of the catheter and to selectively position the distal portion of the catheter during an electrophysiological procedure. The proximal end of the catheter can be manipulated to guide the catheter through a patient's vasculature. The distal tip can be deflected by a pull wire attached at the distal end of the catheter and extending proximally to an actuator in a control handle that controls the application of tension on the pull wire.\nThe foregoing discussion is intended only to illustrate the present field and should not be taken as a disavowal of claim scope."}
{"text": "Ethylenediaminetriacetic acid (ED3A) or its salts (such as ED3ANa.sub.3) has applications in the field of chelating chemistry, and may be used as a starting material in the preparation of strong chelating polymers, oil soluble chelants, surfactants and others. Conventional routes for the synthesis of ethylenediaminetriacetic acid were achieved via its N-benzyl derivative, which was subsequently hydrolyzed in alkaline solutions to ED3ANa.sub.3, thus avoiding cyclization to its 2-oxo-1,4-piperazinediacetic acid (3KP) derivative. Syntheses attempted by both the alkaline condensation of chloroacetic acid with ethylenediamine, and the carboxymethylation of the diamine with formaldehyde and sodium cyanide resulted in complex mixtures requiring complex extraction techniques (e.g. almost exclusive solubility of 3KP in boiling dimethylformamide, Can. J. Chemistry 1970, 48(1), 163-175) to generate the desired product, and then in only relatively poor yield. In addition, conventional processes resulted in large quantities of by-product, such as ethylenediaminetetraacetic acid (ED4A). Where the by-products were especially objectionable, complicated blocking techniques were necessary in order to achieve a relatively pure solution.\nOne example of the synthesis of ethylenediamine-N,N,N'-triacetic acid is shown in Chemical Abstracts 78, Vol. 71, page 451, no. 18369c, 1969. There it is disclosed that ethylenediamine reacts with ClH.sub.2 CCO.sub.2 H in a 1:3 molar ratio in basic solution at 10.degree. C. for 24 hours to form a mixture from which ethylenediamine-N,N,N'-triacetic acid can be separated by complexing the same with Co(III). The resulting cobalt complexes can be isolated through ion exchange.\nThe instant invention is directed to a novel composition of matter that is useful as an intermediate in the synthesis of ethylenediaminetriacetic acid or its salts in high conversions and excellent yield."}
{"text": "In some applications, the resultant thickness of a coating (e.g., a paint) that is applied to a substrate (e.g., the surface of a metal substrate) by a user may be critical, or at least important, to provide desired performance (e.g., proper protection of the substrate). For example, achieving a specified thickness of an applied coating may be critical to preventing corrosion of a metal substrate used in marine applications. Self-inspecting coatings are used in applications such as, for example, marine applications and oil and gas pipeline applications. A self-inspecting coating often includes a coating (e.g. liquid or powder) that provides a visual indication (e.g., visible or invisible to naked eyes) of coating properties (such as thickness). As an example, the visual indication of the coating properties may be provided as the coating is applied or after the coating is applied. For example, a color of the coating can change as the applied thickness changes, in accordance with an embodiment. In this manner, a user is able to perform a certain level of self-inspecting as the user applies the coating. That is, the user may visually observe the color of the coating as it is applied to the substrate in an attempt to determine if the thickness is correct. However, the ability of a user to discern variations in color (and, therefore, variations in the coating film) by observing the coating with the naked eye is limited.\nFurther limitations and disadvantages of conventional, traditional, and proposed approaches will become apparent to one of skill in the art, through comparison of such systems and methods with embodiments of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "Generally, a high voltage semiconductor device may be utilized when high voltage or high current output is required to drive a motor or when a high voltage is input from an external source. Typically, the high voltage semiconductor device includes a high voltage driving region and a low voltage driving region in a system-on-chip structure. In the high voltage device, when a low voltage is applied to a gate electrode and a high voltage is applied only to a drain electrode, the low voltage driving region and the high voltage driving region are formed at the same time. A through gate-oxide implantation (TGI) process is performed during a manufacturing process of a high voltage semiconductor device to form the low voltage driving region and the high voltage driving region on and/or over the chip while maintaining the existing characteristics. In the TGI process, an ion implantation process is performed to form a well region on and/or over a semiconductor substrate on and/or over which a high voltage gate oxide film is deposited.\nAs illustrated in example FIG. 1, a high voltage semiconductor device includes semiconductor substrate 1 having a high voltage device region and a low voltage device region, well region 2, device isolation film 3, gate oxide film 4, gate electrode 5, liner film 6 and interlayer insulating film 7. Device isolation film 3 is formed to define a device isolation region on and/or over semiconductor substrate 1 having the high voltage device region and the low voltage device region. Gate oxide film 4 is formed on and/or over semiconductor substrate 1 in the high voltage device region. After high voltage gate oxide film 4 is formed, a photoresist pattern is formed at a portion of the entire surface of substrate 1. Then, an ion implantation process is performed on semiconductor substrate 1 using the photoresist pattern as a mask to form well region 2 therein. In the ion implantation process for forming well region 2, ions are also implanted into the exposed high voltage gate oxide film 4. Accordingly, a trap site may be formed in gate oxide film 4 into which the ions are implanted. The photoresist pattern is then removed.\nGate electrode 5 is formed on and/or over semiconductor substrate 1 having well region 2. Gate electrode 5 is formed in an active region of substrate 1. Liner film 6 is formed on and/or over the entire surface of semiconductor substrate 1 including gate electrode 5 by forming a tetra ethyl ortho silicate (TEOS) film using preferential metal deposition (PMD). Interlayer insulating film 7 is formed on and/or over the entire surface of the resultant structure including liner film 6. Materials such as hydrogen and boron included in interlayer insulating film 7 may move to the trap site in gate oxide film 4. Consequently, leakage current increases in a threshold voltage region of the high voltage semiconductor device, thereby causing problems of increasing power consumption and reducing the device characteristics. Further, when the TEOS film is used as liner film 6, leakage current of the semiconductor device may increase significantly according to the state and atmosphere of a chamber in which liner film 6 is deposited. Particularly, C3F8 gas is used in cleaning the chamber in which the TEOS film of liner film 6 is deposited. If fluorine included in the C3F8 gas has a high density, a low level of leakage current of the NMOS transistor is measured. On the other hand, if fluorine included in the C3F8 gas has a low density, a high level of leakage current of the NMOS transistor is measured. Meaning, there is a problem that the leakage current of the semiconductor device largely increases according to the fluorine atmosphere in the chamber in which the TEOS film is deposited."}
{"text": "The present invention is particularly directed to elimination of the current practice of using manual labor to remove printed and cut shingled sheets from a web printing press and loading them onto a pallet prior to folding the sheets. Then at a later date, the pallet is transported to an automatic folding machine at which another manual operation is used to remove the sheets from the pallet and to manually jog the sheets into alignment and then to place the sheets into a receiving hopper from which the sheets are fed by a bottom feeder into a folding station at which a device folds the sheets and provides a continuous outflow of printed folded sheets. The cost of these manual operations and the transportation of the pallets adds significantly to the cost of the ultimate printed fold sheet. Also, the pallets take considerable space and the transportation and storage of pallets are also space consuming. Hence, it would be more efficient to eliminate these separate operations and to be able to directly fold the printed sheets issuing from the web printing press.\nIt will be appreciated that the outflow of shingled sheets from the web printing press may be at a high rate of speed, for example, as many as 40,000 sheets per hour and that any folding apparatus connected thereto should operate at a similar high rate of speed without frequent breakdowns so that the web printing press may be kept operating at its full production speed.\nThe sheets leaving the web printing press are shingled on the conveyor and are not accurately aligned on the outcoming conveyor and can not be fed directly to the automatic folding machine because the sheets are too misaligned to be fed automatically from a stack by the existing folding machines. That is, the sheets are often skewed relative to one another with respective sheets being laterally misaligned, or turned slightly, and unevenly spaced in the fore and aft direction on the conveyer with the result that the sheets may not be directly fed into a receiving or storage hopper for automatic refeeding to a folder. Instead, it is the manual operator at the folding machine who corrects the misalignment of the sheets as he lifts the sheets and taps and jogs the same into general alignment prior to placing the aligned sheets into the feed hopper of the automatic folding machine.\nAccordingly, a general object of the present invention is to provide an automatic folding machine with an automatic stacking station at which are collected and aligned incoming sheets so that the sheets may be automatically removed from the stack and folded by the folding machine in a continuous process.\nA further object of the invention is to provide a new and improved automatic stacking and folding apparatus for connection directly to the output conveyor of a web printing press."}
{"text": "1. Field of the Invention\nThe invention relates to a method for fabricating a display, and more particularly to a method for fabricating a flexible display and apparatus and production process thereof.\n2. Description of the Related Art\nFlexible displays have unique advantages such as high impact resistance, light weight and flexibility. As such, in addition to researched applications in newly emerging products such as electronic paper, electronic tags, credit cards, scrolling displays and electronic advertising boards, further applications are being explored for usages in portable electronic products. As for flat panel displays, developmental trends continue to encompass larger areas, lighter weights and thinner frames. For flexible displays, the main developmental trend is for efficient and economic use of a plastic substrate in place of a glass substrate.\nThe conventional flexible display fabrication process, using a plastic substrate, requires steps such as film deposition, photolithography and etching. Also, the apparatuses for manufacturing conventional flexible displays are expensive, and the costs for research and development in this field of technology as well as fabrication are high. Furthermore, the conventional flexible display fabrication process is not a continuous process, thus making it difficult to increase manufacturing yields. As a result, with high costs and high product prices, expanding further application of the conventional flexible displays have been hindered.\nThus, development of a novel method for fabricating a flexible display is desirable."}
{"text": "During drilling operations from just after initial spudding of a well through completion and initiation of production, drilling fluid or drilling “mud” fills the interior of the formed well bore. Some types of muds are petroleum-based materials. Petroleum-based materials comprise at least 90 weight percent of an oil-based mud (OBM) as a continuous phase. Examples of suitable base petroleum materials include crude oil, a distilled fraction of crude oil, including diesel, kerosene, asphalt, waxes, lubricating oils, mineral oil, and heavy petroleum refinery liquid residues. Typically, a minor part of the OBM comprises water or an aqueous solution that is in the mud as an internal phase. Such water-in-oil emulsions are useful to transport chemicals that are not otherwise useful in the continuous phase. Other optional OBM components include emulsifiers, wetting agents and additives that give desirable physical properties to the mud or treat the well bore wall.\nOil-based muds also include synthetic oil-based muds (SOBMs). Synthetic oil-based muds are crude oil derivatives that have been chemically treated, modified, altered or refined, or combinations thereof, to enhance and promote certain chemical or physical properties and exclude other aspects of typical OBMs. SOBMs are monolithic systems that behave in a manner as if they were an OBM but provide a limited and predictable range of chemical and physical behaviors. In comparison to a distilled fraction of crude oil, which may contain several classes (for example, alkanes, aromatics, and heteroatomics) representing thousands of individual compounds, a SOBM usually comprises one class representing at most tens of individual compounds (for example, ester compounds in a C8-C14 range). Examples of useful materials for the base fluid of a SOBM include linear alpha olefins, isomerized olefins, poly alpha olefins, linear alkyl benzenes and vegetable oil- and hydrocarbon-derived ester compounds.\nA mud with an aqueous continuous phase—a water-based mud (WBM)—typically comprises water in a range of greater than 50% to about 99% water. Unlike OBMs, WBMs may have a significant portion of hydrocarbons, including materials that would normally serve as the basis for an OBM, as part of the WBM. The base fluid for the water-based systems include fresh water, natural and saturated salt waters, natural or artificial brines, sea water, mineral water, and other potable and non-potable waters containing one or more dissolved salts or minerals. In regions where water is scarce or environmental regulations do not permit the disposal of untreated formation water, recycling recovered formation water from other production sites can provide an inexpensive source for a WBM, especially if the formation water contains salts and minerals that are useful to stabilize clay and shale downhole.\nBesides salts and minerals, often other additives are useful in attempting to control the viscosity or inhibition of water-based mud. Common additives include sodium or potassium silicates (“silicate muds”) to inhibit shale and seal microfractures that occur during drilling, quebracho (“red mud”) and other tannates, ferrochrome lignosulfonate (“lignosulfate mud”), potassium formate, lignites, phosphates, polyphosphates, gypsum, water-soluble polymers (“polymud”), lime, cellulose and xanthose based polymers, biopolymers, brines, biocides, corrosion inhibitors, foamers and cleaners.\nThe ability to maintain rotational velocity and fluid flow is a significant attribute of all drilling fluids, but this is especially true when the drilling tools stop their rotation and their introduction/withdrawal movement. Fluid momentum and disturbance of the fluid flow within the well bore by the tools permits suspension of solids, incompatible with the continuous phase liquids and gases to be maintained throughout the course of the drilling fluid flow pathway from the surface, downhole, and then back uphole for recovery and reintroduction. Reduction of fluid momentum due not only to general fluid friction but also friction against the sidewalls and the downhole equipment eventually causes the drilling fluid to settle. To prevent this settling, often it is necessary to continue pumping to the surface or slowly rotating the drilling string to keep the drilling fluid moving to a point where solids do not drop out of the continuous phase and incompatible gases and liquids do not separate.\nIt is desirable to include with a drilling fluid a composition that can significantly lower the frictional effects of the drilling fluid such that fluid momentum may be maintained with increased ease by action of the drill string or pumping of the drilling fluid to and from the surface. Such a composition would not only provide safer and more predictable operations with the modified drilling fluid, but also energy usage would be significantly reduced. A composition that is also environmentally friendly and that is biodegradable is also advantageous for use in marine and ecologically-sensitive environments."}
{"text": "1. Field of the Invention\nThis invention relates to emission control of internal combustion engines. In particular, the invention relates to an air/fuel ratio closed loop fuel control of an internal combustion engine equipped with two exhaust gas oxygen (EGO) sensors and a three way catalytic converter. The EGO sensors are located upstream and downstream of the catalyst.\n2. Prior Art\nIt is known that catalyst efficiency is greatly affected by the ratio of air to fuel in the mixture supplied to an engine. If the air/fuel ratio is kept in a narrow range at stoichiometric ratio, catalyst conversion efficiency is high for both oxidation and reduction conversions. Air/fuel stoichiometric ratio is defined as the ratio containing air and fuel in such proportions that in perfect combustion both would be completely consumed, and air/fuel ratio LAMBDA is defined as the amount by weight of air divided by the amount by weight of fuel over air/fuel stoichiometric ratio. The purpose of any closed loop fuel control system is to keep the air/fuel ratio in this narrow range known as a conversion window.\nIt is also known that a control system utilizing one EGO sensor located either before or after a catalyst does not maintain the air/fuel ratio consistently inside the conversion window. Control systems with one EGO sensor located before a catalyst have acceptable time response characteristics but exhibit long term drift due to EGO sensor contamination and aging. On the other hand, control systems with one EGO sensor located after a catalyst have unacceptable time response characteristics but exhibit good long term stability and can indicate a narrow conversion window. Therefore, a control system utilizing advantages of both EGO sensors, i.e., good time response of an upstream EGO sensor and high accuracy of a downstream EGO sensor, is advantageous.\nA number of closed loop fuel control systems utilizing both EGO sensors have been proposed but none are completely satisfactory. U.S. Pat. Nos. 3,939,654 issued to Creps and 4,027,477 issued to Storey describe a dual EGO sensor closed loop fuel control system having two control loops. The first control loop includes an upstream EGO sensor and a proportional or proportional with phase lead controller. The second control loop includes a downstream EGO sensor and a dual integrator controller. This arrangement of the control system precludes the usage of integral or proportional and integral controllers in both control loops simultaneously because such a control system is inherently unstable and can not be made stable by calibration. The drawback of these systems is a low accuracy of the first control loop with a proportional controller. The control accuracy may even become unacceptable when the second control loop is not operational as is the case during initial operation before the downstream EGO sensor reaches its operational temperature.\nOther teachings, represented by U.S. Pat. Nos. 4,831,838 issued to Nagai et al and 4,840,027 issued Okumura et al, utilize a proportional and integral (PI) controller in the first control loop with an upstream EGO sensor. In one embodiment of these patents, calibratable parameters of the PI controller may be modified based on the output of a downstream EGO sensor, the modified parameters being a skip amount, or jumpback, and an integration amount, or ramp. Some other control system parameters such as time delay and reference voltage may also be modified based on the output of a downstream EGO sensor.\nIn another embodiment of the same patents, the output of a downstream EGO sensor is used to generate a second air/fuel ratio correction amount which is used as a multiplier in the main fuel equation (the main fuel equation will be introduced later). In both embodiments, a correction introduced by a downstream EGO sensor has a very low frequency limit cycle superimposed on a relatively high frequency limit cycle produced by an upstream EGO sensor control loop. It results in an undesirable effect known in the control field as beat frequency. Moreover, the initial response of a downstream EGO sensor is so slow that elaborate procedures are incorporated to mitigate this disadvantage. Accordingly, both approaches to a dual EGO fuel control system have been found to be unsatisfactory.\nAnother known control technique using dual EGO sensors, one upstream and one downstream of the catalyst, is a cascade control wherein a signal from the downstream EGO sensor is applied to a summer with a reference signal. The output of the summer is applied to a first proportional and integral controller. A signal from the upstream EGO sensor is applied to a summer and a reference, the reference being the output of the first proportional and integral controller. The output of the second summer is applied to a second proportional and integral controller which then generates the feedback signal to control engine air/fuel ratio.\nIn another scheme similar to the one just mentioned, both of the summers have an applied reference signal. The output of the first proportional and integral controller is not applied to the second summer but instead controls the parameters of the second proportional and integral controller. This is known as parametric control because the parameters of the second controller are controlled by the output of the first controller. Both this system and the previous system are relatively slow in operation. With respect to the later, parametric control, when a parameter is changed, such as the jumpback, ramp, or delay of a control function, it may well take minutes for the effect to be felt. Further, the system is slow to start. It would be desirable to overcome these problems.\nApplicant's invention has a much faster response and uses a single proportional and integral controller having inputs from both the upstream and the downstream EGO sensor."}
{"text": "1. Field of the Invention\nThe present invention relates to a window molding member, for instance for automobiles, as well as a method of manufacturing such molding members.\n2. Description of the Related Art\nAn automobile employs various kinds of molding members, of which a typical example is a window molding member adapted to extend along the periphery of the front or rear window panes, i.e. along a pair of front or rear pillars and the front or rear edge of the roof panel of the automobile body. A variety of requirements are imposed on molding members mainly from design and/or functional viewpoint, and result in an increased demand in the automobile industry for the molding members whose cross-sectional shape varies in the longitudinal direction.\nSpecifically, one proposal is disclosed for example in Japanese Utility Model Application Publication No. 57-54,416, which is directed to a window molding member having an upper portion with a first predetermined cross-section, at least one side portion with a second predetermined cross-section, and at least one corner portion arranged between the side and upper portions. When the molding member is arranged along the periphery of a front window pane, the first cross-section of the upper portion contributes to form a so-called flush outer surface of the automobile body, while the second cross-section of the side portion serves to define at least one channel or weir along the side edges of the window pane. Such an arrangement of the molding member ensures that, during driving in rainy conditions, the weir effectively prevents rain water on the window pane from flowing across the side portion toward the side window, to preserve the driver's and/or front seat passenger's view through the side windows.\nTo produce a window molding member with a cross-section which varies in the longitudinal direction, it is possible to physically divide each molding member into first and second extruded portions with the respectively predetermined cross-sectional shapes, which are connected with each other either by an injection molding process or by using a separate corner connection piece. However, connection of these two portions by means of an injection molding results in formation of undesirable burrs along the junctions and deterioration in the appearance, while use of a corner connection piece results in an increased number of the required components and assembly steps.\nAnother possibility for manufacturing molding members with a longitudinally variable cross-section is disclosed in Japanese Patent Application Laid-open No. 62-231,814, wherein the main body of the molding member has a lip section in the form of a ridge, and two leg sections both to be inserted into a gap between the automobile body panel and the window pane. In use, the first leg section along the upper and side portions of the molding member is engaged by a retainer member on the automobile body panel such that the molding member is retained in position. Furthermore, the second leg section is directly engaged with the periphery of the window plate along the upper portion whereby the lip section is brought into direct contact with the surface of the window plate along the upper portion, while a separate erection member accommodating the second leg section therein is engaged with the periphery of the window pane along the side portion such that the second leg section is supported by the erection member and urges and deforms the lip section away from the surface of the window pane along the side portion to form the weir between the window plate and the ridge. Such an arrangement provides refined appearance due to the continuous and smooth surface along the corner portions integrally connecting the upper and side portions with each other. However, use of erection members along both side portions is sometimes problematic in that, besides an increased number of required components and assembly steps, the lip section tends to be deformed when applied with external force, and it is difficult to stably maintain the desired cross-section of the weir for a long period."}
{"text": "This invention relates to lifting and towing apparatus for vehicles and more particularly to a rigid metal towing bracket adapted to be mounted on a front cross member of a vehicle frame.\nThe prior art is replete with various towing hitches or brackets adapted for mounting on vehicles. An example of one type of hitch bracket is found in the U.S. Pat. No. 4,620,736 issued Nov. 4, 1986, to T. L. Shanks which discloses an adaptor plate for attaching an accessory such as a trailer hitch or winch to the bumper of a vehicle. The plate includes an accessory mount or a wedge insert insertable into a mounting frame or bracket. A tongue and groove assembly cooperates between the bracket and accessory mount for holding the accessory mount in position.\nThe U.S. Pat. No. 2,380,523 issued July 31, 1945 to H. A. Hicks, et al. discloses a vehicle body structure showing a reinforcing member having a pair of longitudinally spaced flanges welded to floor portions of the body. The flanges are stepped to dispose successive flange portions so as to underlie floor portions to which they are respectively welded."}
{"text": "Energy generated from natural resources, such as sunlight and wind, is desirable due to the environmental-friendly factor associated therewith. For example, wind turbines are known for converting rotational movement of the turbine blades, caused by the wind, into electricity. Specifically, the rotation of turbine blades is utilized for the rotation of a shaft, connected to a generator, for producing electricity.\nTo achieve an optimum output of a wind turbine, the rotation of the turbine blades need to be regulated with varying wind speeds/conditions. Specifically, the turbine blades should be allowed to rotate at a set rated revolution per minute (rpm) for achieving the optimum output. However, to maintain the set rated rpm of the turbine blades may be difficult during sudden rise in wind speed or storm. Specifically, rise in wind speed may cause over speeding of the turbine blades, which may cause damage to the wind turbine. Further, fixing such damages may become expensive and time consuming affair."}
{"text": "The present invention relates to an assembly kit including structural elements which are connectable with one another by interengaging undercut grooves and projections. More particularly, it relates to an assembly kit including a plate element and an axle element connectable with one another.\nA plate element has been proposed in the art, having undercut grooves which are adapted to receive undercut projections of other structural elements so as to connect the latter with this plate element. When such several structural elements are connected with one another, there is possibility for insertion of an axle in the plate elements so as to assemble a vehicle structure. It has also been proposed to provide bores in the above plate elements, each of which bores forms an extension of a respective one of the undercut grooves, in order to increase a region of insertion. The axles are inserted into the above bores of the plate element.\nThe above known construction possesses some disadvantages. The positioning of the axle as described above is not suitable for an axle adapted to support wheels of the vehicle structure, inasmuch as the necessary stability thereof is not provided in the region of the undercut groove. Furthermore, at the place where the axle is inserted in the plate element there is no possibility to mount a further structural element on the plate element."}
{"text": "1. Field of the Invention\nThe present invention relates to the FFT (Fast Fourier Transform) used in an OFDM (Orthogonal Frequency Division Multiplexing) system, and more particularly to an FFT method for processing input signals in parallel in order to quickly process the input signals.\n2. Description of the Related Art\nThe basic principle of the OFDM (Orthogonal Frequency Division Multiplexing) system is to convert input data having a high data rate into parallel data which have a low data rate, where the number of parallel data is equal to the number of sub-carriers, and to carry the parallel data on the sub-carriers, respectively, to transmit the data in parallel. The OFDM can reduce relative distortions in the time domain by a multi-path delay spread since the symbol duration of the sub-carrier having the low data rate is increased, and can remove an inter-symbol interference by inserting a protection section that is longer than the delay spread of the channel between OFDM symbols.\nSince the OFDM modulation/demodulation is performed using a plurality of sub-carriers, it is quite difficult to work out its hardware design as the number of sub-carriers is increased. Also, due to the difficulty in keeping the orthogonality between the sub-carriers, it becomes difficult to actually implement the system. Although this problem can be solved by adopting a DFT (Discrete Fourier Transform), the DFT has a drawback in that it requires a large amount of computation. In order to reduce the large amount of computation that is the drawback of the DFT, an FFT (Fast Fourier Transform) has been proposed. Specifically, in the OFDM system, an N-point DFT is required. However, as N increases, the amount of DFT computation also increases in proportion to N2. Accordingly, it is required to provide an algorithm that can efficiently compute the DFT even if N is large. The FFT is an algorithm that remarkably reduces the amount of DFT computation by successively dividing a sequence having a length of N into sequences having a length shorter than N.\nThe FFT of the OFDM performs a computation of a complex number that is composed of a real part and an imaginary part. Accordingly, the real part and the imaginary part are separately inputted by hardware, and in designing a processor that performs the FFT, an inverse FFT (IFFT) can be performed by changing the positions of the real part and the imaginary part with each other. The FFT may be implemented in an array type or in a pipeline type. The array FFT structure is very complicated and enlarged by hardware, and thus its implementation is almost impossible if the number of FFT computation points is large. By contrast, the pipeline FFT structure is regular, is relatively easy to control and makes a serial input/output possible, and thus it is most frequently used in application fields that require a high performance.\nHereinafter, the DFT and the FFT will be explained in order. Signals having a predetermined period which are expressed by the DFT are defined by Equation (1):\n                              X          ⁡                      (            k            )                          =                              ∑                          n              =              O                                      N              -              1                                ⁢                                          ⁢                                    x              ⁡                              (                n                )                                      ⁢                          ⅇ                              j                ⁢                                                      2                    ⁢                                                                                  ⁢                    π                                    N                                ⁢                nk                                                                        (        1        )            \nwherein N denotes the number of signals, k denotes 0 to N−1, x(n) denotes an input signal and X(k) denotes an output signal. As described in Equation (1), the amount of DFT computation is increased as the value of N is increased.\nFIG. 1 is a view exemplifying a process of performing a conventional FFT. Hereinafter, an algorithm that performs the conventional FFT will be explained in detail with reference to FIG. 1.\nIn particular, FIG. 1 describes a case in which the point (N) of the FFT is 16. Signals inputted to the FFT are x[0] to x[15]. Hereinafter, for the convenience in explanation, horizontal lines from x[n] to X[n] are called computation lines. Referring to FIG. 1, the FFT is composed of first to 16th computation lines. x[0] is divided at point a and transferred to the first computation line and the ninth computation line at point b. x[1] is divided at the point a and transferred to the second computation line and the tenth computation line at the point b. x[14] is divided at the point a and transferred to the seventh computation line and the 15th computation line at the point b. x[15] is divided at the point a and transferred to the eighth computation line and the 16th computation line at the point b.\nThe first to eighth computation lines at the point b add the transferred signals and output the added signals, and the ninth to 16th computation lines at the point b subtract the transferred signals and output the subtracted signals. The signals outputted from the point b are transferred to point c. The computation lines at the point c perform the same operations as the computation lines at the point a.\nThe first to fourth computation lines and the ninth to 12th computation lines at point d add the transferred signals and output the added signals, and the fifth to eighth computation lines and the 13th to 16th computation lines at the point d subtract the transferred signals and output the subtracted signals. The signals outputted from the point d are transferred to point e. The computation lines at the point e perform the same operations as the computation lines at the point a. The first to second computation lines, the fifth to sixth computation lines, the ninth to tenth computation lines and the 13th to 14th computation lines at point f add the transferred signals and output the added signals. The third to fourth computation lines, the seventh to eighth computation lines, the 11th to 12th computation lines and the 15th to 16th computation lines at the point f subtract the transferred signals and output the subtracted signals. The signals outputted from the point d are transferred to the point e.\nThe signals outputted from the point f are transferred to point g. The computation lines at the point g perform the same operations as the computation lines at the point a. The odd-numbered computation lines at point h add the transferred signals and output the added signals, and the even-numbered computation lines at the point h subtract the transferred signals and output the subtracted signals. Through the above-described process, the FFT is performed with respect to the input signals.\nHowever, the FFT has the problems in that as the computation points N are increased, it takes a lot of time to process the input signals. This is because the FFT as illustrated in FIG. 1 does not process all the signals in parallel, but processes only a part of the signals in parallel. Accordingly, the necessity of an FFT that can quickly process the input signals is on the rise."}
{"text": "These inventions relate to methods, an apparatus for their implementation, of unique player participation games, and for improved methods of play for games of chance. More particularly, these inventions relate to new and improved games involving player participation in a broadcast medium, such as television, and in other communication media, such as over the Internet or other communications network.\nPlayer participation games fall broadly under the categories of games of chance and games of skill. One of the main forms of games of chance is lotteries, which by definition, involve the three elements of: 1) prize, 2) chance and 3) consideration. If these three elements are present, then the game is considered to be a lottery, and is typically then run by a governmental entity. In the United States, lotteries are typically run by the individual states, or collectively by a group of states. In other countries, it is typically the national government that runs the lottery. Countries and states attempt to strictly limit the game play to their geographic boundaries. For example, in Austria, while electronic access to the game may be available over the Internet, or in order to play, the person must have a bank account in Austria, and be able to navigate the non-english menu.\nGames have been conducted in any of a number of formats. Certainly, live, in person games have been performed. Yet other games have been played and broadcast over a broadcast medium, such as radio or television. Yet other games have been played through active communication media, such as the telephone, or over a communication network such as the Internet.\nVarious attempts have been made to provide game play over the Internet. By way of example, the game show Jeopardy has been placed on the web at http://www.sony.com.\nVarious other attempts have been made to extend the general concept of gambling to broad communication media, such as the Internet. For example, U.S. Pat. No. 5,800,268 entitled, xe2x80x9cMethod of Participating in a Live Casino Game from a Remote Locationxe2x80x9d has been asserted in a litigation in against an off shore corporation. The \"\"268 patent discloses a system in which a player may participate in a live casino game from a location remote from the casino. A player interface station, such as a computer terminal or other special input device, is connected by a communication line to the casino. A second communication line is established from the casino to the player\"\"s financial institution. The player is presented with an image of an actual xe2x80x9clivexe2x80x9d game. The player then participates directly as if they were physically present at the casino. A wager is cleared with the player\"\"s financial institution to insure adequate resources to cover the bet.\nU.S. Pat. No. 4,845,739 to Ronald A. Katz is entitled, xe2x80x9cTelephonic Interface Statistical Analysis Systemxe2x80x9d. The patent describes various operating formats, including a format the to be performed in association with television media. Specifically, in one embodiment, a real-time format is provided in which television viewers participate on a real-time basis in a game show for prizes. Expanded audience participation is achieved. Various levels of qualification are provided, such as for a child\"\"s television game format is utilized, parental clearance may be required. The use of personal identification numbers (pin numbers) is disclosed. In one implementation, the caller is prompted to identify which of the actual studio of audience participants the caller will be aligned with. Additionally, the caller may be instructed to indicate the extent of a wager. As the game progresses, the individual player\"\"s accounts are credited or debited, thereby providing on-going accounting data. In yet another implementation, a non real-time operation is provided. Such a show might involve a quiz for callers based on their ability to perceive and remember occurrences within the show. Pre-registration is optionally utilized. In this implementation, a sequence or time clock would be utilized in order to limit or control individual interfaces to a specific time or geographic xe2x80x9cwindowxe2x80x9d. In this way, the caller questions may be utilized across various time zones without the caller having obtained the question earlier than other callers within a given time zone.\nBerman, U.S. Pat. No. 5,108,115 discloses a game show and method entitled xe2x80x9cInteractive Game Show and Method for Achieving Interactive Communication Therewithxe2x80x9d. An interactive communication system is provided which permits individuals to electronically select at least one possible outcome of a plurality of outcomes of a future event. Successful contestants possibly share in a prize aware associated with the event. A home audience of a televised game show may electronically communicate a series of random numbers using their touch tone telephone to participate in the show.\nRecently, various governmental entities and trade organization have addressed the issue of game play over the Internet. Congressman Kye has introduced a bill which would preclude the offering of Internet based gaming, though permitting states to offer Internet gambling. Consideration has been given to requiring that the states sponsored gaming be limited to an intranet, in an effort to limit those participating to persons physically resident within the states boundaries. Various international lottery organizations have promoted similar restrictions, namely, precluding the individuals offering of games of chance, and reserving that option exclusively to the state.\nVarious lottery formats are known to the art. In one classic format, a predetermined number of tickets are provided with certain printed matter, such as numbers or other indicia, where the information is then obscured by a scratch off layer. By removing the layer and revealing the underlying information, the ticket holder may determine whether they have won or not. Various extensions have been made to a xe2x80x9cvirtualxe2x80x9d scratch off ticket where no physical is provided.\nA conventional lottery proceeds as follows. First, a series of numbers are selected, either by the player or by some automated selection system, such as by computer. Upon the occurrence of a pre-determined event, such as on a set date and time, numbers are randomly chosen. Both mechanical methods, such as selection of ping-pong balls bearing numeric designations, or electronic means such as through a random number generator, may be utilized. The selected numbers are then provided to the participants, such as through a broadcast medium like newspapers, radio and television. Finally, the holder or holders of winning the tickets then present their ticket for payment.\nIn yet another aspect of game play, a typical television presented game show lasts on the order of one half hour. Various shorter format games or shows have been utilized, for example, a football based advertisement or game has been presented by IBM during televised football games under the name xe2x80x9cyou make the callxe2x80x9d. Yet other shorter version games have been presented over web TV or on the game show network.\nThe television game show xe2x80x9cWho Wants to be a Millionairexe2x80x9d is believed to have originated in Britain, and has become extremely popular in the United States. The game is a trivia game. While being principally a game of skill, the nature of the questions, or the contestants knowledge of the potential answers, makes the game at times a guessing game or game of chance. The format consists of one contestant and one host. The contestant is presented with a question and four possible answers. If the contestant answers the question correctly, they advance to a next level, each level being associated with a higher monetary prize amount, which is roughly twice the amount of the preceding level. A contestant is given three xe2x80x9clife linesxe2x80x9d: a xe2x80x9c50/50xe2x80x9d where in two incorrect answers are removed, thereby leaving the correct answer and one incorrect answer, the xe2x80x9cphone a friendxe2x80x9d, wherein the contestant may call a friend by telephone and solicit their response to the question, subject to a 30 second time limit, and an xe2x80x9cask the audiencexe2x80x9d option where the audience is polled regarding their view of the correct answer to the question. Various safe levels are established, such as at $1,000.00 such that the contestant would be awarded that amount of money in the even that they fail to correctly answer a question. Finally, after a question is posed, the contestant may elect to discontinue play, and to receive that amount of money won at the preceding level.\nDespite the wide spread participation in various forms of game play, as well as the suggestions for implementing those games on a mass communication network, such as through the telephone or Internet, the possibility for new games, or improved game play exists. In particular, there is a need for improved games of chance, which provide excitement for the player, and optionally a viewer audience.\nThis invention relates to methods and associated apparatus for novel game play. In the preferred embodiment, the game is a game of chance.\nIn the preferred embodiment, the game is played at a multiple number of levels. At each level, the contestant is presented with multiple options, such as a depiction of four uniquely labeled boxes, amongst which the contestant may choose. The options would include at least one positive outcome and at least one negative outcome. In the case of four boxes, e.g., one could include a strike, two could include a monetary amount, which may be either the same or different and optionally, the fourth box could comprise a mystery box, described below. The contestant selects, at random, one of the options. If the option selected is one of the positive options, such as a monetary amount, they proceed to the next level and the winnings are added to the prior winnings total. If a negative option is selected, such as a strike, in the preferred embodiment, the level is reset and play continues at that level. Preferably, the player is allowed a predetermined number of negative events, such as three strikes, prior to discontinuing play.\nThe xe2x80x98mystery boxxe2x80x99 consists of a decision within a decision. A first decision was to select that option, which then was revealed as comprising a mystery box. The player is then given the option of whether to reveal that option. The option within the mystery box would include at least one positive result and at least one negative result. In the preferred embodiment, there would be three results possible with a mystery box, a positive result such as a multiplier for the money, such as a doubler of the contestantxe2x80x99 prior winnings, an updating of the safe level for the player or an additional monetary amount. Alternatively, other positive results such as a free play or a reduction in the number of negative events is possible. Preferably, the probability of a negative result from the opening of the mystery box should be equal to the probability of a negative event if the mystery box were not selected.\nThe prizes at the various levels may be set as desired to result in a predetermined pay out for the game. Optionally, guaranteed low end prize structures (GLEPS) may require payment of predetermined prize amounts, and possibly payment of a minimum amount of a prize e.g., $500.00. The monetary spacing between various levels may be set as desired, either as an arithmetic progression or as a multiplicative progression, e.g., a substantial doubling of the prize amount at every level. Optionally, when a maximum game level is reached a jackpot or other proportionally large prize may be awarded. If the jackpot is not won in a given game, it may then roll over to a subsequent game. Alternative forms of progressive play may be utilized.\nIn another aspect of this invention, game play in a first game may require progression through a plurality of levels, leading to game play on a second game for those who have reached the maximum level on the first game. In one implementation, the maximum prize level in the first game may be equal to the minimum prize level in the second game.\nVarious modes of play are contemplated. In studio game play may be utilized with a broadcast, either live or for taped replay. Yet another mode of game play involves playing at a gaming venue, such as where other games of chance, e.g., slot machines, are played. Yet another venue may consist of game play by the player from their hotel room in a venue which allows gambling. In yet another mode of game play, a network, such as the internet, may be utilized to permit game play, whether for a monetary amount or to provide other points or indications of score. The game may be played in any venue where not prohibited, whether on land or in an airplane or ship, and may be played in any form of wired or wireless environment, such as via hand-held web enabled communication devices.\nThe game may be played by a single individual, or may be played with multiple players. The multiple players may play against one another, for scoring, or may merely play in parallel without further interaction.\nAccordingly, it is an object of this invention to provide an improved game of chance having a higher level of audience interest and potential participation.\nIt is yet another object of this invention to provide for an improved Internet game of chance.\nIt is yet a further object of this invention to provide for enhanced modes of game play in association with existing forms of game play."}
{"text": "Various machines, such as mining trucks, are known to employ drive propulsion systems to propel or retard the machine, such as a mechanical drive or an electric drive. An electric drive propulsion system, for example, generally includes an alternator, or other electrical power generator, driven by an internal combustion engine. The alternator, in turn, supplies electrical power to one or more electric drive motors connected to wheels of the machine. The motors are generally connected to the wheels by way of a final drive assembly that reduces the rotational speed of the motor. The final drive of a machine may employ one or more planetary gear sets to reduce the output speed of the propulsion system. Such planetary gear sets may be partially submerged in an oil bath for cooling and/or lubrication. At high rotational speeds, turbulence and churning created in such oil baths may contribute to the loss of energy and/or the generation of heat.\nSome gear train assemblies attempt to reduce turbulence and churning by providing a stationary shroud that fits closely to the gears and fully encloses the sides and outer diameters of meshing gears. For example, U.S. Pat. No. 5,048,370 (the “'370 Patent”) teaches a shroud for enclosing several gears within a gear train. The shroud taught in the '370 Patent includes an input and output port for injection and ejection, respectively, of cooling fluid. However, a close fitting, stationary shroud may not be an effective or feasible solution for reducing turbulence created by large planetary gears that are used in the powertrain of many machines, particularly planetary gear sets having rotating carriers that are partially submerged in an oil bath."}
{"text": "Problems associated with fuel lubricity arose in the mid- 1960's when a number of aviation fuel pump failures occurred. After considerable research, it was realized that advances in the refining of aviation turbine fuel had resulted in the almost complete removal of the naturally occurring lubricating components from the fuel. The removal of these natural lubricants resulted in the seizure of fuel pump parts. By the mid-1980's, it seemed likely that a similar problem was imminent in diesel fuel pumps. Fuel injection pump pressures had been steadily increasing while there was also a growing concern to reduce the sulfur content of the diesel fuel. The desire to reduce the sulfur content of the diesel fuel, in an effort to reduce pollution, required the use of more rigorous fuel refining processes. It was determined that as refining processes became more stringent, the naturally occurring oxygen containing compounds and polyaromatics which contribute to diesel fuel's inherent lubricity were eliminated.\nEnvironmental concerns have led to a need for fuels with reduced sulfur content, especially diesel fuels. However, the refining processes that are used to produce fuels with low sulfur contents also result in a product of lower viscosity and a lower content of other components in the fuel that contribute to its lubricity, for example, polycyclic aromatics and polar compounds. Furthermore, sulfur containing compounds in general are regarded as providing anti-wear properties and a result of the reduction in their proportions, together with a reduction in proportions of other components providing lubricity, has been an increase in reported failures of fuel pumps in diesel engines using low sulfur fuels.\nThis problem may be expected to become worse in the future because in order to meet stricter requirements on exhaust emissions, high-pressure fuel pumps are being introduced and are expected to have more stringent lubricity requirements than present equipment.\nIn certain types of in-line diesel injection pumps, engine oil contacts diesel fuel. Engine oil may also come into contact with the diesel fuel through direct addition of used engine oil to the fuel. Certain types of lubricity additives used in low sulfur diesel fuel have been found to contribute to fuel filter blockage and to pump plunger sticking. Lubricity additives having poor compatibility with engine oil have been shown to cause these problems. Compatibility is defined as the tendency for the diesel fuel containing the lubricity additive not to form fuel insoluble deposits, gels or heavy sticky residues when in contact with engine oil. These deposits, gels or residues have been shown to cause fuel filter blockage and injection pump sticking. The additives of the present invention are compatible with engine oil.\nMannich reaction products have been taught for use as detergent/dispersants in fuels, primarily gasoline, for years. The prior art Mannich reaction products typically contain high molecular weight alkyl groups on the hydroxyaromatic compounds. In contrast, the Mannich reaction products of the present invention are obtained from alkyl-substituted hydroxyaromatic compounds wherein the alkyl group contains from 9 to 30 carbon atoms.\nU.S. Pat. No. 3,877,889 discloses Mannich bases useful as additives for liquid fuels to impart dispersancy, anti-icing and rust inhibiting properties. The reference fails to teach the use of said Mannich reaction products as lubricity additives in low sulfur compression ignition fuels.\nU.S. Pat. No. 4,231,759 teaches reaction products obtained from the Mannich condensation of high molecular weight alkyl-substituted hydroxy aromatic compounds, amines and aldehydes for improving the detergency of liquid hydrocarbon fuels.\nU.S. Pat. No. 5,853,436 discloses diesel fuel compositions containing a lubricity enhancing amount of a salt of an alkyl hydroxyaromatic compound and an aliphatic amine. These salts are different than the reaction products of the present invention.\nWhile the prior art is replete with numerous treatments for fuels, it does not disclose the addition of the present additives to low sulfur compression ignition fuels or teach their use for providing enhanced lubricity to said fuels."}
{"text": "The present invention relates to a braid and a braiding machine, and more particularly, to a polygonal braid and a braiding machine therefor, in which a guide plate having a track capable of braiding a polygonal braid and a feed gear corresponding thereto are provided. The polygonal braid braided in a square or triangular shape can be utilized as a binding string for shoes or clothes etc because of increased binding force by a polygonal. Braids of various quality and colors may be created by using different qualities or colors of strands.\nBraids are utilized in several fields, for example, as part of an electric wire or hose, as a binding string etc. A specific braid is formed on the outer circumference of the electric wire or the string and provides an elastic and relaxed covering for an interior electrc wire or a string etc. and protects the interior electric wire from being contaminated or damaged by impact, braids are often used in place of string for daily use in shoes or clothes etc., in addition to specialized uses.\nA general braider is composed of a guide plate having a track on which a spool is moved, a feed gear for moving an electric spool along the track the guide plate, a driving gear for driving a plurality of feed gears and a plural number of rollers on which a braided wire is wound, etc.\nFIG. 6a shows a guide plate for manufacturing a general cylindric braid and its braid. On this guide plate 100, two tracks 101, 101xe2x80x2 on a gentle circular curve line of a jig jag shape are formed, intersected with each other. As shown in FIG. 6b, on its lower part, a plurality of feed gears 102, 102xe2x80x2 are positioned beneath and aligned within the intersected curves of tracks 101, 101xe2x80x2. In such construction, the plurality of feed gears 102, 102xe2x80x2 are simultaneously rotated by a rotation of the driving gear 103, yarn from separate spools are combined within feed gears 102, 102xe2x80x2 onto the guid plate 100 which is rotated.\nTherefore, a plural number of spools of yarn are rotated, repeatedly performing a rotation and a revolution centering around a center point of the guide plate 100, feeding out yarn, which are intersected with one another, rotating along the track 101, 101xe2x80x2.\nOn an outer circumference of the central yarn, thus, a braid based on a cylindric shape is produced by the rotation of the spool as shown in FIG. 6c. \nThe ordinary cylindrical braid as described above, when used as a binding of shoes, has a low frictional force due to the cylindrical shape which can result in the shoe lace coming loose.\nWhan a braid is made using a single color, by prior art methods, the brain color is monotonous; when using a single color, by prior art methods, the braid color is monotonous; when using various colored braid, the braid color may appear untidy.\nIn order to overcome the problems presented in prior art braiding methods, the present invention teaches a braid formed in a polygonal shape such as a triangle or a square shape with each edge of each face the polygon intersecting with an edge of the adjoining face of the polygon."}
{"text": "Some conventional vehicles include an electrical wiring harness extending between a roof and a roof liner of the vehicle. The electrical wiring harness is attached to the roof liner using tape and also using two-piece clips. A first piece of each clip is glued to the roof liner and a second piece of each clip is taped to the electrical wiring harness and is snapped into the respective first piece."}
{"text": "1. Field of the Invention\nThe present invention relates to encoding of data to provide for robust error recovery due to data losses typically incurred during transmission or storage of signals.\n2. Art Background\nA number of techniques exist for reconstructing lost/damaged data due to random errors that may occur during signal transmission or storage. However, these techniques are limited at handling the loss of consecutive packets of data. Consecutive loss of packets of data is described in the art as burst error. Burst errors may result in a reconstructed signal with such a degraded quality that it is easily apparent to the end user.\nAdditionally, compression methodologies used to facilitate high speed communications compound the signal degradation caused by burst errors, thus adding to the degradation of the reconstructed signal. Examples of burst error loss affecting transmitted and/or stored signals is seen in high definition television (xe2x80x9cHDTVxe2x80x9d) signals, mobile telecommunication applications, as well as video storage technologies including compact disc (CD), video disk (e.g., DVD), and video cassette recorders (VCRs).\nIn one application, the advent of HDTV has led to television systems with a much higher resolution than the current standards proposed by the National Television Systems Committee (xe2x80x9cNTSCxe2x80x9d). Proposed HDTV signals are predominantly digital. Accordingly, when a color television signal is converted for digital use the luminance and chrominance signals can be digitized using eight bits. Digital transmission of NTSC color television so encoded requires a nominal bit rate of about two hundred and sixteen megabits per second. The transmission rate is greater for HDTV which would nominally require about 1200 megabits per second. Such high transmission rates are well beyond the bandwidths supported by current wireless standards. Accordingly, an efficient compression methodology is required.\nCompression methodologies also play an important role in mobile telecommunication applications. Typically, packets of data are communicated between remote terminals in mobile telecommunication applications. The limited number of transmission channels in mobile communications requires an effective compression methodology prior to the transmission of packets. A number of compression techniques are available to facilitate high transmission rates.\nAdaptive Dynamic Range Coding (xe2x80x9cADRCxe2x80x9d) and Discrete Cosine Transform (xe2x80x9cDCTxe2x80x9d) coding provide image compression techniques known in the art. Both techniques take advantage of the local correlation within an image to achieve a high compression ratio. However, an efficient compression algorithm can result in compounded error propagation because errors in an encoded signal are more prominent when subsequently decoded. This error multiplication can result in a degraded video image that is readily apparent to the user.\nThe present invention provides a method for compressing data by determining a central value that is greater than the minimum value and less than the maximum value of the range of data. In one embodiment, the central value is chosen to be a value that substantially reduces a decoding error in the event that the range of values is subsequently estimated. In one embodiment, the central value is the value that minimizes the expected mean square error during reconstruction when there is an error. In one embodiment, the maximum and minimum values represent intensity data for pixels of an image. In another embodiment, the compression process is Adaptive Dynamic Range Coding, and the central value is a value within the dynamic range, excluding the maximum and minimum values."}
{"text": "It is well recognized in the petroleum industry that boron containing compounds are desirable additives for lubricating oils. One such boron containing compound is disclosed in U.S. Pat. No. 3,224,971 to Knowles et al. which relates to intracomplexed borate esters and to lubricating compositions containing said esters. The borate esters are organo-boron compounds derived from boric acid and a bis (o-hydroxy-alkylphenyl) amine or sulfide.\nAnother extreme pressure lubrication composition is disclosed in U.S. Pat. No. 3,185,644 to Knowles et al., which relates to lubricating compositions containing amine salts of boron-containing compounds. The amine salts are formed by reaction of a hydroxy substituted amine and a trihydrocarbyl borate. The amine-borate compounds thus formed are described as useful as load carrying additives for mineral and synthetic base lubricating oils.\nBoric-acid-alkylolamine reaction products and lubricating oils containing the same are disclosed in U.S. Pat. No. 3,227,739 to Versteeg. These amine type products are prepared by reacting equal molar proportions of diethanolamine or dipropanolamine and a long chain, 1, 2-epoxide. The intermediate reaction product thus produced is reacted with boric acid to produce the final reaction product. These compounds are added to lubricants to prevent rust formation.\nAnother boron ester composition is described in U.S. Pat. No. 3,269,853 to English et al. which discloses a boron ester curing agent which consists of a cyclic ring structure containing boron, oxygen, nitrogen, carbon and hydrogen.\nAnother boron composition is disclosed in U.S. Pat. No. 3,598,855 to Cyba which relates to cyclic borates of polymeric alkanolamines formed by reacting a borylating agent with a polymeric alkanolamine. The compounds thus formed are described as additives for a wide variety of petroleum products including lubricating oils.\nCurrently, there are phosphorus-containing additives which provide extreme pressure, anti-wear and/or friction-reducing properties to automotive engine oils. However, with the advent of the catalytic converter, alternative additives are needed. During combustion in an automotive engine, any oil which leaks or seeps into the combustion chamber yields phosphorus deposits which poison the catalyst in the catalytic converter. As a result, there is a need for automotive engine oil additives which are phosphorus-free but provide useful extreme pressure, anti-wear, and/or friction-reducing properties to the oil.\nAccordingly, it is one object of the invention to provide a phosphorus-free additive having such properties and which, upon combustion, will not adversely affect the catalyst in the automotive catalytic converter.\nIt is yet another object of the present invention to provide boron-containing, heterocyclic compounds or derivatives thereof which have extreme pressure, anti-wear and friction-reducing properties.\nYet another object of the present invention is to provide a lubricating composition having extreme pressure, anti-wear and friction-reducing properties.\nA further object of the present invention is to provide a lubricating composition containing extreme pressure, anti-wear, friction-reducing and corrosion prevention additives, and in addition, an anti-oxidant to prevent attack of oxidants upon metal bearings.\nOther objects and advantages of the invention will be apparent from the following description."}
{"text": "Spring retractable rule assemblies have been available commercially for many years. One of the most desirable characteristics commercial rule assemblies can possess is a relatively long blade standout. To date, as a practical matter, that standout length of most blades has rarely exceeded 7 feet or at most approximately 9 feet. Standout is generally measured by the length of the rule assembly blade that can be extended in a self-sustaining manner without the blade buckling under its own weight. An important characteristic of standout relates to the vertical bend that the blade takes in its maximum self-sustained extension. This is generally expressed in terms of the height the housing assumes from a horizontal surface when the free end of the blade is just touching the horizontal surface. The vertical height of the housing above the horizontal surface, the position at which the free end of the blade touches the horizontal surface and the vertical projection of the position of the housing onto the horizontal surface roughly define three points of a right angle triangle. The hypotenuse of the triangle represents a close approximation of the actual length of the blade extending from the housing and the horizontal leg of the triangle represents the linear horizontal extension of the blade. It is generally recognized to be desirable to maintain the ratio between the linear horizontal extension to the actual extension as near to one as possible. There always exists a need to provide a retractable rule assembly that will provide greater standout with a greater ratio of linear horizontal standout to actual standout."}
{"text": "The microcomputer revolution began two decades ago, but for most of that time it was required only that the system being designed comply with rather unforgiving requirements of size (on the order of a desktop) and power consumption (several dozens or hundreds of watts). The early personal computers used large numbers of discrete components, but thereafter it became commonplace to use \"chip sets\" which reduce the computer system to half a dozen integrated circuits each with dozens or hundreds of pins, or preassembled \"mother boards\", either of which leaves little or no room for optimization by the individual system designer. Thus the individual designer cannot do much in the way of reducing power consumption or changing physical size or form factor.\nIn more recent years, however, the marketplace has come to demand computer systems, such as personal computer systems, which run independently of AC (mains) power and which are meant to be carried from place to place and used in portable fashion. In such systems there is a renewed attention to issues of power consumption, weight, and size. One consequence of the greater attention to power consumption is the development of communications channels and protocols according to which system elements which provide and consume power are in communication with each other to permit sophisticated power management. It is desirable that the power management communications channel and protocol satisfy several requirements, for example, small pin count (so that batteries need not have too many connector pins) and undemanding protocols (so that devices can be slow if necessary). One approach is to employ a synchronous bus, that is, a bus in which data is passed with reference to a clock line. The clock line is \"high\" in a quiescent state, and is pulled low if a device on the bus (the \"bus master\") wishes to pass a bit of data on the line for reading by any or all of the devices on the bus. At a later time, the bus master raises the clock line and again pulls the line low to indicate that a subsequent bit of data is readable by any or all of the devices on the bus. In this way a message composed of many data bits is communicated across the bus.\nTo accommodate a range of types of bus devices with varying response times and latencies, it is desirable to define a \"clock stretching\" element of the protocol. According to this aspect of the protocol, any bus device, having noted that the clock line has been pulled low, can itself pull the clock line low. Indeed in the general case it is assumed that any number of bus devices may have done so. During the time that the clock line remains low, the defined behavior of the bus master is to maintain the data level on the data line. In this way, a bus device can take as long as desired to read the data value (and to prepare for the reading of subsequent data values). In colloquial terms, it can be said that the clock line remains low until the slowest of the bus devices has managed to get up to speed and to process data on the data line.\nThose skilled in the art will appreciate that power consumption in a microprocessor or microcontroller is monotonically (and generally linearly) related to the clock speed thereof. Thus the system designer who is attempting to maximize battery life (or to minimize power consumption) will consider a variety of measures including switching a microcontroller to a very slow clock speed, or indeed powering down the microcontroller, during times of low or zero workload. (This may be termed \"putting the controller to sleep\".) For example the designer of a microcontroller for a computer keyboard may actually power down the microcontroller except when a key has been pressed. If the user makes a thousand keystrokes, the microcontroller is powered up and down a thousand times.\nReturning now to the above-mentioned synchronous bus with clock stretching, it may happen that the system designer chooses to have a bus device go to sleep, only to be awakened when there is activity on the bus. But a block of data transmitted on the bus may be intended for the very device that is asleep, and yet it is desired that no data be lost. The \"clock stretching\" aspect of the protocol may be employed to prevent such data loss. The circuitry that accomplishes the \"clock stretching\" cannot itself be put to sleep, of course, but must be kept functioning at all times in case data is transmitted on the bus. The clock stretching circuit cannot have any clocks running continuously. In at least one known prior art design, it is necessary that the clock stretching circuitry be served by a clock that runs continuously, at a megahertz or so. This leads to non-negligible power consumption.\nA further concern is that there be minimal power leakage into any particular bus device while it is asleep.\nFIG. 10 shows a typical prior art circuit that permits microcontroller 8051SL to go to sleep if desired and includes a clock stretching function. Four flip-flops (\"start det\", \"clk hold\", \"busy\", and \"stop\") are required along with several logic gates and comparators. There are many components to be assembled during manufacture and they take up space and consume power.\nIt would be extremely desirable to provide a \"clock stretching\" circuit that would permit putting a microcontroller to sleep so as to save power, such a circuit having substantially smaller component count, power consumption, and assembly cost. It would be additionally desirable to provide a level shifting circuit that isolates the bus device from the bus when the bus device is in a power saving mode, has its power removed, or is operating at a power supply voltage lower than that of the bus itself."}
{"text": "Developmental disorders such as autism spectrum disorders (ASD) affect nearly 14% of children in the United States. Diagnostic methods for conditions such as ASD vary considerably, and even the use of “best practice” tools provides rather poor sensitivity and specificity to the conditions. Late diagnosis of developmental disabilities reduces effectiveness of treatments and often results in poor outcomes. Furthermore, treatment providers (e.g., pediatricians or other medical professionals) lack adequate tools for measuring progress in these conditions."}
{"text": "Cytotoxic T lymphocytes (CTL) and natural killer (NK) cells perform tumour surveillance and provide a defense against viral infection and intracellular pathogens, by inducing apoptosis of virus-infected or transformed cells. A major component of this defense is the glycoprotein perforin. Upon stable conjugation of the CTL or NK cell with a target cell, perforin is released, binds calcium and assembles into aggregates of 12-18 molecules that form trans-membrane pores in the plasma membrane. This allows leakage of cell contents and the entry of secreted serine proteases (granzymes) which promote apoptosis.\nStimulation of CTL and NK cells, leading to abnormal cellular destruction, occurs in several autoimmune diseases (e.g., insulin-dependent diabetes) and in therapy-induced conditions (e.g., allograft rejection, graft-versus-host disease). In this context, small-molecule inhibitors of perforin function are of potential interest as a new class of therapeutic immunosuppressive agents.\nTo date, the only reported direct inhibitors of perforin function are those published by the present inventors (Lena et al, J. Med. Chem., 51(23), 7614-7624, 2008; Lyons et al, Bioorganic & Medicinal Chemistry, 19, 4091-4100, 2011; Spicer et al, Bioorganic & Medicinal Chemistry, 20, 1319-1336, 2012). Other reported inhibitors of perforin function are non-selective, complex natural products, primarily concanamycin A and other V-ATPase inhibitors such as bafilomycin A and prodigiosin 25-Cs that inhibit acidification. Other reported non-selective perforin inhibitors include cytochalasin D (an inhibitor of actin polymerisation), antimycin A and oligomycin A (inhibitors of cell respiration) and some protein kinase inhibitors (calphostin C, herbimycin A, staurosporine). However, such non-selective compounds display a broad spectrum of biological effects that generally make them undesirable for use in the treatment or prevention of conditions associated with aberrant perforin expression and/or activity.\nIn one or more aspects, the present invention may advantageously provide a class of compounds and their analogues as drugs for immunosuppressive therapies, or to at least provide a useful alternative to existing treatment modalities."}
{"text": "Online media services allow users to access their media on a wide range of devices, including devices owned by other people. For example, when a friend is hosting a guest, the guest may use the host's home media player to access and play music over the internet from the guest's online media account. In many circumstances a guest will remain logged in to the host's media player after accessing a particular song because the guest and friend do not want to go through the log in process again in order to have access to additional music, or because the host's media account does not contain the same type of music. This can result in the guest forgetting she was logged in or getting distracted and not paying attention to other people who may be accessing the media player. Once the guest logs in to the host's media player, any person with access to the media player may access aspects of the guest's media account unless the guest logs out. For example, the guest's purchasing account may be accessed to purchase additional media without the guest's consent, or the guest's credentials may be used to view the guest's confidential data or gain access to other online accounts belonging to the guest. Thus, conventional techniques for allowing users access to their online media content often place users at risk for significant financial, privacy, and data security concerns."}
{"text": "1. Field of Invention\nThis invention relates generally to cigarette lighters having a child-resistent mechanism and more specifically to lighters employing a double-button child-resistent mechanism.\n2. Related art\nCigarette lighters containing piezoelectric units are very useful and have become quite prevalent in modern times. Cigarette lighters of the type described herein generally contain a lighter housing that is small enough to be held in the palm of an adult hand. The operation of piezoelectric cigarette lighters is somewhat simpler than that of the traditional flint/spark-wheel lighter. Generally, the lighter is operated by depressing an actuator button, which both activates the piezoelectric unit and acts on a fuel-release lever to release fuel. As a result, a flame is produced at a location opposite the actuator button. As is evident, this process avoids the need for operation of a spark wheel simultaneously with operation of a fuel-release button in order to generate a flame. Obviously, there is an advantage to the simplicity that is offered by piezoelectric cigarette lighters. On the other hand, in the hands of children, or others who do not know how to safely and properly operate the lighter, such lighters are as dangerous as any other spark and/or flame-producing device. Therefore, a need has been realized to equip cigarette lighters with safety features that minimize accidental or improper use by inexperienced persons, especially young children.\nMany inventions have been created to address this safety-related concern. Generally, these inventions have sought to introduce safety mechanisms that disable operation of the actuator button of the lighter. As such, these lighters normally consist of a safety feature whereby the operational path of the actuator button is blocked by a latch, button, slide, or other blocking means. Proper operation of the lighter requires that the blocking means be moved out of the path of the actuator button, or other structure that might be integral with the actuator button, before a flame can be produced. Only then is the operator able to depress the actuator button and produce a flame. As such, the prior art requires additional structural members, as well as additional steps (e.g., lateral or longitudinal disengagement of a blocking means), to operate the lighter.\nIn some of the aforementioned cigarette lighters, the safety mechanism is passive. That is, once the safety feature is de-activated by moving the blocking member from the \"locked\" to the \"unlocked\" position, the lighter remains in the \"unlocked\" position, and thus is operable as a cigarette lighter with no safety feature at all. In these devices, the lighter remains in the \"unlocked\" position until the safety feature is activated again by manually re-engaging the safety mechanism (e.g., by manually returning the blocking means to the \"locked\" position).\nIn order to address this problem, some inventions have introduced safety mechanisms that are activated automatically after each use of the lighter. In general, this improvement has alleviated some of the fears associated with leaving the lighter in an \"unlocked\", operable position after the operator has finished using the lighter. Nevertheless, a disadvantage that is common to the passive, as well as the active, cigarette lighters is that their operation is usually quite cumbersome. Frequently, in order to use such cigarette lighters, the operator must use more than one finger, and sometimes more than one hand, to perform several functions simultaneously. As such, loss of ease of use is the price that is paid for any additional amount of safety that might be achieved.\nTherefore, there is a need for a device that not only achieves the stated safety goals, but also is amenable to operation with relative ease. The invention described herein offers such a combination and consists of a safety button that is similar in size and physical location to the conventional activation button. The invention requires that an ignition button, located in a cavity within the safety button, be depressed simultaneously with the safety button before a flame can be produced. In this way, young children are coaxed into believing that they can operate the lighter in the usual way, i.e., by pressing only the safety button. However, such operation will produce neither a spark nor a flame. Moreover, given the relatively small size of the ignition button, operation of this button requires an amount of strength and pulp that are rarely found in the fingers of young children. At the same time, due to the placement of the ignition button, simultaneous operation of both the safety button and the ignition button requires use of only one finger, so that operation of the lighter by the intended adult user is no different from operation of a lighter with no safety mechanism at all."}
{"text": "Makeup cases of this type generally contain a grille over a shallow cup of makeup product of powder cake, as well as an applicator puff enclosed between the grille and the cover. These cases can only be used dry.\nThere are also known cases of the same type containing a moistenable sponge for moist use."}
{"text": "1. Field of the Invention\nThe present invention relates to a connector assembly that engages a second connector with a first connector by inserting a second housing of the second connector into a first housing of the first connector and connects a terminal of the second housing to a terminal in the first housing.\n2. Description of the Related Art\nThere has conventionally been known a connector assembly that engages a second connector with a first connector by inserting a second housing of the second connector into a first housing of the first connector and that connects a terminal of the second housing to a terminal in the first housing. A connector assembly that, by inserting a housing of a male connector into a housing of a female connector, mutually connects terminals of both the housings is an outstanding example thereof.\nThis kind of connector assembly is used, for example, to connect various types of electric cables with each other inside a vehicle. Because electric cables are often routed within a confined space in a vehicle, various schemes are incorporated into the connector assembly to fix the connector assembly on the vehicle side in a narrow small space of the vehicle. As an example, there is a connector assembly that is fixed on a vehicle side by forming a cassette type insertion section on a side face of a housing of a female connector and/or a male connector and inserting an end of a bracket fixed on the vehicle side into the cassette type insertion section for clamping (for example, refer to FIG. 5 in Patent Literature 1: Japanese Patent Application Laid-Open Publication No. H08 (1996)-330026).\nThe connector assembly described above, in some cases, is formed with a high-voltage specification with recent widespread use of electric vehicles (EVs) and hybrid electric vehicles (HEVs). The connector assembly with such a specification has a tendency of being further upsized than a conventional connector assembly with a low voltage specification. Thus, the connector assembly tends to have large vibrations thereof when vehicle vibrations propagate through a bracket. Accordingly, in a large-sized connector assembly with a high voltage specification, such as a shield connector, it is preferable that both female and male connectors are tightly coupled to each other without any looseness under an engagement state thereof.\nIn this view, a conventional connector assembly which is lighter and smaller than a connector assembly with a high voltage specification has a sufficient coupling force in an engagement direction sufficiently ensured to completely connect terminals mutually. However, because of no necessity of such a force in each direction within a plane intersecting with the engagement direction, the connector assembly is not structured to obtain a sufficient coupling force with an engagement force between a female connector and a male connector. Accordingly, a connector assembly with a high voltage specification having a high tendency of being upsized may require a new structure to provide enhancement of a coupling force, which a conventional connector assembly has not adopted."}
{"text": "Injection of a liquid such as a drug into a human patient or an agriculture animal is performed in a number of ways. One of the easiest methods for drug delivery is through the skin which is the outermost protective layer of the body. It is composed of the epidermis, including the stratum corneum, the stratum granulosum, the stratum spinosum, and the stratum basale, and the dermis, containing, among other things, the capillary layer. The stratum corneum is a tough, scaly layer made of dead cell tissue. It extends around 10-20 microns from the skin surface and has no blood supply. Because of the density of this layer of cells, moving compounds across the skin, either into or out of the body, can be very difficult.\nThe current technology for delivering local pharmaceuticals through the skin includes methods that use needles or other skin piercing devices. Invasive procedures, such as use of needles or lances, effectively overcome the barrier function of the stratum corneum. However, these methods suffer from several major disadvantages: local skin damage, bleeding, and risk of infection at the injection site, and creation of contaminated needles or lances that must be disposed of. Further, when these devices are used to inject drugs in agriculture animals, the needles break off from time to time and remain embedded in the animal.\nThus, it would be advantageous to be able to inject small, precise volumes of pharmaceuticals quickly through the skin without the potential of a needle breaking off in the animal."}
{"text": "Drill rods or pipes come in sections that are joined by male threads, tapered convergently, and female sockets, threaded and tapered internally complementarily to the male threads. The drill strings are rotated in a direction that tends to tighten the joints. That, plus the fact that the joints are liable to get gritty material between the threads, makes breaking the joints difficult. Various power mechanisms have been proposed and used (e.g., U.S. Pat. Nos. 4,345,493, 5,727,432). Where non-rotating vises have been employed, they have held the rod while a pipe wrench or the like has been used to turn the section that is not gripped by the vise. The wrench has been turned either manually or by some power mechanism. The pipe wrench type joint breakers have the disadvantage that the gripping surface is limited, putting extreme pressure in a limited area, and may not provide sufficient gripping force. Another approach to turning the free section has been to provide jaws sliding in a cradle rotated on a radius concentric with the radius of a pipe or rod clamped between the jaws. However, this construction has heretofore been subject to wear, with metal-to-metal contact, and a tendency to distortion, which leads to misalignment between the sections of the string.\nOne of the objects of this invention is to provide an improved open top rotating vise, in which there is substantially no wear or distortion, and in which the rotation of the rotating vice is facilitated.\nOther objects will become apparent to those skilled in the art in the light of the following description and accompanying drawings."}
{"text": "Social media is pervasive in today's society. Friends keep in contact throughout the day on social networks. Fans can follow their favorite celebrities and interact on blogs, micro-blogs, and the like. Such media are referred to as “social media,” which can be considered media primarily, but not exclusively, for social interaction, and which can use highly accessible and scalable communication techniques. Brands and products mentioned on such sites can reflect customers' interests and feedback.\nSome technologies have been developed to analyze social media. For example, some systems allow users to discover their “influence scores” on various social media. An influence score is a metric to measure a user's impact in social media."}
{"text": "Portable computing devices are increasingly powerful. Likewise more progress is being made towards paperless environments. Portable computing devices can support the goals of achieving paperless environments."}
{"text": "1. Field of the Invention\nThe present invention relates to a controller for image processing apparatus. More specifically, the present invention relates to a controller which is utilized in an image processing apparatus such as a television game machine wherein it is required for an operator or player to continuously send quick and adequate reactions or responses to the television game machine while the operator or the player is watching a television screen, and outputs an electrical signal for controlling an image on the television screen according to a lean angle when an operating portion of the controller is leaned by hands of the operator or the player.\n2. Description of the Prior Art\nAs a technique for a controller device which can be utilized for a television game machine and utilizes a lean or orientation of an operating portion, for example, a Japanese Patent Application Laying-open No. 58-112576 (corresponding to U.S. Pat. No. 4,503,299) for \"Control-Lever for a Game\" (hereinafter, called as \"first prior art\") is known. In the control-lever for game, a plurality of movable balls are contained in the operating portion, and when the operating portion is leaned, the balls roll toward a leaned direction, and therefore, switches provided on inner side walls of the operating portion are turned-on or -off, whereby presence or absence of the lean or orientation can be detected. Furthermore, the first prior art also discloses a technique in which mercury is used instead of the balls.\nOther than the above described prior art, a technique in which the presence or absence of the lean as well as a degree of lean are somewhat considered as disclosed in \"Freestanding Multidirectional Electrical Control Device\" of U.S. Pat. No. 4,445,011 (hereinafter, called as \"second prior art\"). The second prior art utilizes mercury or electrical conductive liquid as similar to the first prior art, and the degree of lean is detected by combining contact electrodes made of material having an electrical resistance value which is changed in accordance with its length and the mercury or the electrical conductive liquid.\nOn the other hand, as an input device for moving a cursor on a screen of a computer video monitor, which is not for the television game machine, \"Computer Input Device Using an Orientation Sensor\" of U.S. Pat. No. 5,068,645 (hereinafter, called as \"third prior art\") is known. This device is one of a few prior arts that the device is mounted on a head of an operator such that an orientation angle of the head can be continuously detected, and used as a computer input device. The third prior art is an orientation angle detecting technique wherein a liquid is half-filled in a spherical housing, and a change of a transmission light which is outputted as a result of refraction at a boundary between the liquid and an air half-filled in an upper portion of the housing when the liquid is oriented is utilized.\nThe first prior art detects only whether the operating portion is leaned or not, and can not detect the lean angle. Furthermore, there is a problem in safety that the mercury is utilized for consumer products such as television game machines.\nFurthermore, the second prior art utilizes the mercury or electrical conductive liquid as similar to the first prior art, and therefore, the second prior art can not be utilized for a purpose that a delicate movement of the operator is continuously detected within a wider angle range in view of utility, reliability, and operability.\nThe third prior art has a limit of detectable orientation angle of approximately .+-.30 degrees because of structural feature of the orientation angle detector, and therefore, the third prior art is limited to an input device which is mounted on the head of the operator. A use environment of the input device is also limited to a computer environment wherein the operator can use the input device quietly or gently because of a vibration of a surface of the liquid, responsivity and etc. If the third prior art is utilized in an environment for a television game machine wherein it is possible to presume that the operator or player moves lively, since the surface of the liquid waves, a reflection state of the light is undesirably changed, and accordingly, it is possible to consider that it is difficult to put the third prior art into practical use."}
{"text": "1. Field of the Invention\nThis invention relates to a method and an apparatus for truing a grinding wheel, and more particularly, to a method and an apparatus for truing a grinding wheel having a non-straight cylindrical grinding surface.\n2. Prior Art of the Invention\nIn a numerical controlled grinding machine, a grinding wheel having a stepped peripheral shape is used for grinding a workpiece having a stepped outer surface, namely, a workpiece having plural cylindrical outer surfaces and plural shoulder end surfaces adjacent to the respective cylindrical outer surfaces. Although such grinding wheel is manufactured to have a desired stepped peripheral shape, the grinding wheel has an initial shape slightly different from the desired shape because of inaccuracy of the manufacturing process. Further, there is a case where a user wants to slightly change the peripheral shape of the grinding wheel before using the grinding wheel.\nAccordingly, it is necessary to true the grinding wheels before using them, so that the grinding wheels have desired peripheral shapes.\nAn apparatus which can be used for the above-mentioned truing operation is disclosed in the U.S. Pat. No. 4,899,718 which was assigned to the assignee of this application. The apparatus is provided with a contact detection sensor such as an AE (acoustic emission) sensor for outputting a contact signal when a truing tool contacts a grinding wheel. In dressing operation, the truing tool is moved relative to the grinding wheel so that the truing tool contacts each of the two inclined surface portions of the grinding wheel, and then the position of the truing tool is detected when the truing tool contact the respective surface portions. Based on these positions thus detected, it is judged which surface portion has a larger removal amount to be removed by truing operation, compared to the other. An initial truing start position of the truing tool is determined based on the radial position of the surface portion which has a larger removal amount, and the truing tool is then traversed along the outer surface of the grinding wheel while following a template. Ascertainment is made as to whether the contact sensor outputs a contact signal continuously during the traverse fed. If the continuous issue of the contact signal is not ascertained, then the truing tool is infed a predetermined amount against the grinding wheel before it is traversed. The infeed and the traverse feed of the truing tool are repeated until the continuous issuance of the contact signal during each traverse feed movement is finally ascertained.\nThe above-mentioned truing apparatus has a disadvantage that the mechanical structure thereof is complex, because it needs a template and a stylus, and that the template must be changed when the grinding wheel is trued in different shapes.\nFurther, in the conventional apparatus, it is required that the AE sensor outputs the contact signal continuously during the time period when the truing tool contacts the grinding wheel. However, it is difficult to obtain a stable contact signal from the AE sensor, because the AE sensor detects a sound wave which is generated when the truing tool contacts the grinding wheel. The level of the sound wave changes depending on the shape of a surface portion which the truing tool contacts, namely, depending on whether the surface portion has a straight shape or a curved shape. Further, the level of the sound wave changes depending on the amount of coolant supplied to the outer surface of the grinding wheel during truing operation. Therefore, it is difficult to judge whether or not the truing tool continuously contacts the grinding wheel, even though the detection sensitivity of the AE sensor is adjusted. This causes inaccuracy of truing operation."}
{"text": "1. Field of the Invention\nThe present invention relates to demodulators specifically, to a demodulator capable of avoiding phase and amplitude shifts between a first modulated signal and a second modulated signal by making the input impedance of a first matching section employed in a splitting/matching section for receiving the first modulated signal equal to the input impedance of a second matching section which is also employed in the splitting/matching section for receiving the second modulated signal.\n2. Description of the Related Art\nIn recent years, high-speed transmission technologies making use of a high-frequency band such as a millimeter-wave frequency band are being intensively and extensively researched with an aim to transmit a signal at smaller power consumption and a lower transmission cost due to a smaller circuit scale through the use of CMOS (complementary metal oxide semiconductor) technology. A signal transmission apparatus making use of a high-frequency band is configured to employ a modulator for transmitting a modulated signal of a millimeter-wave frequency band and a demodulator for receiving the modulated signal from the modulator and generating the modulated signal.\nA demodulator 600 employed in the existing signal transmission apparatus is explained. As shown in FIG. 10, a modulated signal received by an antenna 510 is amplified by an amplifier 520. The amplified signal is split into a first modulated signal and a second modulated signal at a branch point Bo which is provided between the amplifier 520 and a squaring circuit 530. The squaring circuit 530 is a section for multiplying the first modulated signal by the second modulated signal in order to demodulate the modulated signal. That is to say, the squaring circuit 530 generates a demodulated signal as a result of the multiplication. The squaring circuit 530 supplies the demodulated signal to an amplifier 540 which then amplifies the demodulated signal and outputs the amplified signal.\nIn addition, the demodulator also referred to as a signal detection circuit can have another proposed typical configuration described as follows. In this proposed configuration, for example, a signal output by the signal detection circuit is compared with a reference voltage and a direct-current component of the result of the comparison is supplied to an IF (intermediate frequency) amplifier and fed back to the signal detection circuit. See, Japanese Patent Laid-open No. Sho 57-37905 (hereinafter, as Patent Document 1). In accordance with this signal detection circuit, a direct-current voltage of the detection output can be made stable."}
{"text": "A gas turbine engine includes a compressor for pressurizing ambient air which is then mixed with fuel in a combustor and ignited for generating combustion gases which flow through a turbine which extracts energy therefrom. The turbine includes stator vanes which preferentially channel the combustion gases through a row of rotor blades which in turn rotate a rotor disk for providing shaft power. Since the combustion gases are hot, the stator vanes and the rotor blades are typically internally cooled using a portion of the compressed air bled from the compressor.\nMore specifically, the vanes and blades typically include a hollow airfoil having an internal cooling flow channel through which is channeled the cooling air bled from the compressor for internally cooling the airfoil. Convective heat transfer cooling is typically enhanced by providing a variety of techniques, including tubulators within the airfoil which have various conventional forms. The cooling air may simply be channeled through the airfoils, or the airfoils may include trailing edge discharge apertures or film cooling holes along either the pressure or suction sides of the airfoil, or both, in accordance with conventional practice. These outlets discharge the cooling air from the airfoil directly into the combustion gases and are suitably sized to provide a minimum backflow pressure margin to prevent the combustion gases from flowing into the airfoils through these outlets.\nAny compressed air which is not used for generating combustion gases, such as that used for cooling the airfoils, decreases the overall operating efficiency of the engine and, therefore, increases specific fuel consumption (SFC). Accordingly, airfoils are cooled by a myriad of different arrangements which attempt to maximize the cooling thereof while minimizing the amount of cooling air bled from the compressor. This must also be done while maintaining acceptable backflow margin and without undesirable pressure losses which decrease the efficiency thereof.\nOne advanced concept for cooling gas turbine engine airfoils simulates the airfoils in a laboratory environment and uses a rotating chopper to provide pulsed or intermittent flow for convection cooling. The simulation demonstrates that convective heat transfer coefficients may be increased using pulsed flow over steady flow for constant airflow; however, the mechanical assembly of moving parts used in the test are clearly impractical for use in a gas turbine engine such as one used for powering an aircraft in flight."}
{"text": "1. Field of the Invention\nThe present invention relates to a current sensor that detects a current flowing in a to-be-measured current path, and more particularly to a current sensor that detects a current flowing in a to-be-measured current path by using magneto-electric conversion elements.\n2. Description of the Related Art\nTo control or monitor an electronic device, a well-known current sensor is attached to a to-be-measured current path to detect a current flowing in the current path. As this type of current sensor, a current sensor that uses magneto-electric conversion elements such as Hall elements or magneto-resistive elements is known. It is known that a plurality of elements are used to improve the sensitivity of the magneto-electric conversion elements and reduce adverse effects by external magnetic fields.\nA current sensor 900, which is an proposed conventional current sensor, has a concave cut into which a to-be-measured current path (not illustrated) is interposed, as illustrated in FIG. 17 (see Japanese Unexamined Patent Application Publication No. 2001-066327). FIG. 17 is a perspective view that schematically illustrates the structure of the current sensor 900 described in Japanese Unexamined Patent Application Publication No. 2001-066327. The current sensor 900 in FIG. 17 has a case 920 that has a concave cut 911 in which a to-be-measured current path is interposed, a circuit board 910 having a cut 921, the circuit board 910 being placed in the case 920, and magneto-electric conversion elements (here, Hall elements) 930 and 931 that output electronic signals according to the intensity of a magnetic filed generated by a current flowing in a conductor placed in the vicinity of the concave cut 911. Thus, the current sensor 900 that can be made compact and easy-to-integrate can be provided.\nAnother known current sensor has four magnetic impedance elements that are oppositely placed, centered around an electric wire holder, which holds an electric wire (see EP1037056 A1). This current sensor has an arc-shaped opening, which works as the electric wire holder, at an engagement portion of a case having a convex part and a concave part. On the concave part side of the case, the four magnetic impedance elements are oppositely placed around the opening in its circumferential direction so as to be equally spaced.\nTo achieve higher measurement precision than in the above conventional technologies described in Japanese Unexamined Patent Application Publication No. 2001-066327 and EP1037056 A1, current sensors having more magneto-electric conversion elements can be contemplated. FIGS. 13A and 13B are plan views of current sensors in comparative examples. FIG. 13A illustrates a current sensor according to a first comparative example, in which eight magneto-electric conversion elements C15 are positioned at positions at which they surround a to-be-measured current path CB. FIG. 13B illustrates a current sensor according to a second comparative example, in which six magneto-electric conversion elements C25 are positioned at positions at which they surround the to-be-measured current path CB. To simplify an explanation, only the to-be-measured current path CB, neighboring current paths CN, and the magneto-electric conversion elements C15 and C25 are illustrated.\nAs illustrated in FIGS. 13A and 13B, the current sensors in the comparative examples each have a plurality of magneto-electric conversion elements (C15 or C25) positioned on a circumference centered on the center of the to-be-measured current path CB in a plan view; each two adjacent magneto-electric conversion elements (C15 or C25) are spaced at the same angular interval. Specifically, the magneto-electric conversion elements C15 of the current sensor in the first comparative example in FIG. 13A are placed at the vertexes of a regular octagon. The magneto-electric conversion elements C25 of the current sensor in the second comparative example in FIG. 13B are placed at the vertexes of a regular hexagon. Thus, even if the position of the to-be-measured current path CB is shifted a little, when detected values from the magneto-electric conversion elements (C15 and C25) are added, it is possible to prevent measurement precision from being easily lowered."}
{"text": "Karst topography is distributed in soluble rock areas all over the world, and has a total area of 51×106 square kilometers, accounting for about 10% of total area of the earth. A karst distribution area in China is about 0.91-1.3 million square kilometers, wherein karst mountain regions in the Yunnan-Guizhou Plateau under subtropical monsoon climate have the worst environment and the most prominent contradiction between people and land. Due to an unscientific production development manner of local people and factors such as centralized monsoon climate rainfall, huge mountain elevation differences and the like, as well as a particularly obvious water and soil loss phenomenon caused by unique duality and vulnerability of the karst environment, productivity and living standards of people have been severely influenced. How to simulate a karst rock desertification environment under different conditions and obtain accurate water and soil loss data to serve karst study, in order to help people understand the karst environment and an underground water and soil loss formation mechanism and to guide people to realize development according to local conditions in a special environment has become urgent problems to be solved."}
{"text": "1. Field of the Invention\nThe present invention relates to a recyclable formwork whose used condition can be easily known. The recyclable formwork is represented by a plastic formwork. The recyclable formwork, however, may be not only the plastic formwork but also a metal formwork, especially a formwork of aluminium, iron, or steel, as long as characteristics thereof, such as weight, price, or the like are suitable.\n2. Description of the Related Art\nCountermeasures against greenhouse gas represented by carbon dioxide are important subjects in which human beings should immediately aim to solve. Carbon dioxide reduction should be strongly carried out without damaging the economy. The inventors will discuss a new proposition in a field of formworks used for curing concrete.\nUntil now, formworks are mainly concrete panels made of tropical wood.\nInstead of the concrete panels of wood, plastic formworks are developed and used. Some types of the plastic formworks can be perfectly recycled and are regarded as next-generation formworks.\nThere are, however, the following problems with respect to plastic formworks. First, a unit price of the plastic formworks is about five times that of concrete panels. Secondly, initial investment for the plastic formworks is extensive. Furthermore, the investment cannot be amortized in one construction site.\nAs for costs considering a number of repeated use times, the first number of repeated use times of the concrete panels are normally about five, and the second number of repeated use times of the plastic formworks is about fifty, that is, ten times the first number. As for costs per one time use, which is obtained by dividing the unit prices by the numbers of repeated use times, the concrete panels are almost double the plastic formworks. The plastic formworks are superior to the concrete panels with respect to not only environmental load but also the cost per one time use.\nHowever, since there is no means for effectively managing the long use time of the plastic formworks, the plastic formworks have not been practically and widely used yet.\nThe formworks are brought to a construction site to be built there. Next, concrete is cast between the built formworks to be cured, the concrete is hardened and solidified, and then the formworks are released. However, curing time remarkably differs according to portions of a structure to be built, such as a wall, a beam, and a floor, or the like.\nThere is no technical means for managing the formworks taking the differences of curing time caused by the portions where the formworks are used into consideration. Actually, a construction supervisor of the construction site and the men use the formworks on their hunches.\nIn short, it is impossible to solve the above problem with conventional plastic formworks and it is necessary to add a new function to the conventional plastic formworks. Furthermore, consideration must be done such that the added new function does not spoil the other functions of the plastic formworks. A technique that satisfies such requirement has not been known yet.\n[Document 1] Japanese patent application Laid-open on No. 2004-332301"}
{"text": "Continuously running equipment is often encountered in which signals are produced in synchronism with its operation. For example, demodulation signals are often needed for modulating systems, especially for multiple-beam spectral photometers, regardless of whether the modulating systems are driven by synchronous motors, as they generally have been (e.g., DD-PS No. 65 468) or by stepping motors (e.g., U.S. Pat. No. 4,386,852, DD-PS No. 228 058). In U.S. Pat. No. 4,386,852, a higher-frequency timing signal is disclosed for the purpose of producing therefrom, by means of a divider circuit, the control frequency of the stepping drive, and also for the purpose of forming a demodulation signal by means of an additional divider synchronized by the optical signal by means of a null detector.\nThis system, which is usable only for single-beam apparatus, also has the important disadvantage that, in the case of higher absorptions of the optical sample, the demodulation is impaired by noise.\nIn DD-PS 228 058, the pulses of the stepping drive control frequency are fed to the stepping drive, and are also carried through an adjustable delay circuit followed by pulse shortening. They control a counter logic such that the demodulation signals, which serve for the analog processing of the measurement, will be in phase with the electrical analog signals formed from the optical signals.\nAn optocoupler cooperating with the modulating mirror supplies for this purpose a synchronization signal which at the beginning of the modulation period resets the counter logic. This eliminates the disadvantages of U.S. Pat. No. 4,386,852, but since adjustment can be made with the adjustable delay circuit only within one step, it is necessary to pre-adjust the modulating mirror and optocoupler with a stepping motor energized in a defined manner, or to assemble them together in an apparatus with sufficient accuracy.\nIn addition to the complexity involved in the construction of the apparatus, the necessary pre-adjustment interferes with service to the modulating system including the optocoupler.\nAlso, in the case of aging phenomena, phasing errors occur between optical signals and analysis signals as the load torque changes at the stepping motor, and necessitate manual readjustment.\nIt has been proposed (GDR Pat. No. 242,089), for synchronization with respect to a given required-phasing signal, to make the stepping motor run one step per modulation cycle slower or faster by means of control frequencies differing from the synchronous frequency, until the required-phasing signal and sensor signal are sufficiently in agreement. Here the same disadvantages are encountered as in DD-PS No. 228 058.\nIt is furthermore generally known to operate a stepping drive not by pulses of a control frequency, but for example to produce commutator signals by means of a computer, whereby the windings of the stepping motor are switched."}
{"text": "The use of portable electronic devices has increased significantly in recent years, with many applications typically residing in the memory of such devices. Exemplary applications include messaging applications, calendar applications and social media applications. Electronic devices often receive communications for these applications, which contain information of importance to users. These electronic devices then often provide notifications that correspond to the received communications.\nExemplary communications include instant messages, calendar invitations, social media updates, microblog posts and news stories. Exemplary notifications associated with these communications may include digital images, video, text, icons, control elements (such as buttons) and/or other graphics to notify users of the receipt of these communications. Exemplary applications receiving communications and generating notifications include instant messaging applications (e.g., iMessage from Apple Inc. of Cupertino, Calif.), calendar applications (e.g., iCal or Calendar from Apple Inc. of Cupertino, Calif.), social networking applications, microblogging applications, and news applications.\nBut user interfaces for accessing notifications, and methods of navigating to and from such interfaces, are cumbersome and inefficient. For example, the notifications may be displayed in a confusing manner, and navigation to and from interfaces that contain notifications may also be confusing. These methods take longer than necessary, thereby wasting energy. This latter consideration is particularly important in battery-operated devices."}
{"text": "1. Field of the Invention\nIn at least one aspect, the present invention relates to methods of coating a substrate with a clearcoat and compositions thereof and, in particular, to methods of coating a substrate with a clearcoat by applying to the substrate a dual curable composition that is first photocurable and then thermally cured.\n2. Background Art\nTypically, the painted surfaces of an automobile are protected by coating with a clearcoat. Clearcoats protect the vehicle from the deleterious effects of sunlight. Accordingly, these coatings typically have light stabilizers, usually consisting of a combination of UV absorbers and free radical scavengers. The absorbers prevent the energetic rays of the sun from causing permanent damage to the polymer matrix of the clearcoat and the underlying coats, including pigments. The free radical scavengers deactivate the highly reactive species that arise as a result of unwanted breakdown processes, and act to promote further breakdown.\nCurrently, there are two main categories of clearcoat compositions that are used to form clear coatings. These categories are medium solid coating systems and high solid coating systems. Solid as used in this context refers to components that are not volatile organic compounds (VOC) including liquids with low vapor pressure. Medium solid coating systems typically contain volatile organic solvents in amounts over 70 weight percent. Accordingly, these systems are undesirable because of environmental and health concern. Moreover, such high solvent systems are subject to government regulations in many countries. High solid systems are more desirable because such systems contain much less volatile organic solvents. In the high solid systems, solvents are typically replaced by liquid oligomers or liquid monomers.\nAlthough high solid coating systems are desirable because of the relatively low amounts of VOCs, coatings from such systems often produce coatings marred by sag. Sag refers to the phenomenon of runs and drips that occurs in paint coating. The tendency of a coating to sag results from several factors. For example, sag may occur from edge effects generated from localized high film build around edges, holes in the substrate, character lines, and the like. Sag may also result from the increased surface tension due to solvent evaporation on two surfaces at an edge and by Faraday's Cage effect. Moreover, sag is observed to be thickness dependent. In the absence of flow control agents, reducing the film thickness by a factor of two reduces the sag by a factor of four. Additionally, for a coating containing 3% microgel a similar reduction in coating thickness results in a 12-fold reduction in sag.\nCoating reaction kinetics is another factor that needs to be considered in minimizing sag. The maximum temperature reached by a coating prior to gel is an important parameter because it essentially determines the minimum viscosity of the coating after solvent evaporation. Therefore, with respect to sag, it is desirable to utilize coatings with higher rates of reaction. These systems become cross-linked sooner, building molecular weight which increases until the coating gels, thereby avoiding or limiting sag. The temperature at gel is higher for high solids coatings versus medium solids coatings because of the higher extent of reaction at gel for high solids coatings. Accordingly, high solid systems have an increased sag potential due to this phenomenon.\nCoating viscosity and cure conditions are additional important factors in controlling sag. If the viscosity is high just after and during cure, then sag may be avoided. However, for high solids coating systems, the low molecular weight resin typically used in these systems and the extent of cure at gel makes sag somewhat inevitable. A significant difference in fluidity between medium solids (“conventional”) clearcoat and high solids clearcoat has been observed. Specifically, during a thermal cure cycle the medium solids clear maintained limited fluidity over the range of heating rates, whereas the high solids clear are significantly affected by the heating rate (lower heating rates—resulting in greater fluidity). Similarly, higher molecular weight systems produce limited fluidity as compared to lower molecular weight systems.\nTypically, in medium solid coating systems, sag is minimized by the use of large amounts of VOCs during the application and cure of the coating. That is, the high molecular weight resin used in these systems require large amounts of organic solvent(s) to reduce the high molecular weight resins viscosity within the wet coating. High solid coatings use lower molecular weight resins to bring down the viscosity. In doing so, thermal cure sag tolerance has been compromised. In high solids coating systems, reduction of sag depends on rheological control agents (“RCA”) to modify the flow and deformation of the liquid coating system. Key characteristics that are sought when adding Theological control agents are to limit settling within the coating, to improve atomization by shear thinning during spray application, and to avoid sag during the thermal cure cycle of the coating by quickly reestablishing a high viscosity after application.\nAccordingly, there exists a need from an improved clearcoat composition that contains low amount of volatile organic solvents and produces a coating with low sag."}
{"text": "A Multiple-Input Multiple-Output (MIMO) has recently been in the spotlight in order to maximize the performance and communication capacity of a wireless communication system. MIMO technology is a method which breaks away from technology using one transmission antenna and one reception antenna and can improve the transmission efficiency of transmission/reception data by adopting multiple transmission antennas and multiple reception antennas. An MIMO system is also called a multiple antenna system. MIMO technology is the application of technology for gathering and completing data pieces received by several antennas without being dependent on a single antenna path in order to receive one entire message. Consequently, the data transfer rate may be improved in a specific range or the range of a system for a specific data transfer rate may be increased.\nMIMO technology includes a transmission diversity, spatial multiplexing, beamforming, and so on. The transmission diversity is technology for sending the same data through multiple transmission antennas in order to increase transmission reliability. The spatial multiplexing is technology for sending data at high speed without increasing the bandwidth of a system by sending different data through multiple transmission antennas at the same time. The beamforming is used to increase the Signal to Interference plus Noise Ratio (SINR) of a signal by applying a weight according to a channel state in multiple antennas. The weight may be represented by a weight vector or a weight matrix and called a precoding vector or a precoding matrix.\nThe spatial multiplexing includes spatial multiplexing for a single user and spatial multiplexing for multiple users. The spatial multiplexing for a single user is called a Single User MIMO (SU-MIMO), and the spatial multiplexing for multiple users is called Spatial Division Multiple Access (SDMA) or a Multi-User MIMO (MU-MIMO). The capacity of an MIMO channel is increased in proportion to the number of antennas. An MIMO channel may be divided into independent channels. Assuming that the number of transmission antennas is Nt and the number of reception antennas is Nr, the number of independent channels Ni is Ni≦min{Nt, Nr}. Each of the independent channels may be said to be a spatial layer. A rank is the number of non-zero eigenvalues of an MIMO channel matrix and may be defined as the number of spatial streams that can be multiplexed.\nMIMO technology includes a codebook-based precoding scheme. The codebook-based precoding scheme is a method of selecting a precoding matrix which is the most similar to an MIMO channel, from among predetermined precoding matrices, and sending a Precoding Matrix Index (PMI). This method can reduce the overhead of feedback data. A codebook consists of codebook sets which can represent spatial channels. In order to increase the data transfer rate, the number of antennas has to be increased. With an increase of the number of antennas, a codebook has to consist of a more number of codebook sets. If the number of codebook sets increases according to the increased number of antennas, not only the overhead of feedback data may be increased, but also there is a difficulty in designing the codebook.\nThere is a need for a method to which a codebook-based precoding scheme can be efficiently applied in a multiple antenna system having an increased number of antennas as compared with the existing number of antennas."}
{"text": "Molecular weight distribution (MWD), or polydispersity, is a well known variable in polymers. The molecular weight distribution, sometimes described as the ratio of weight average molecular weight (M.sub.w) to number average molecular weight (M.sub.n) (i.e., M.sub.w /M.sub.n) can be measured directly, e.g., by gel permeation chromatography techniques, or more routinely, by measuring I.sub.10 /I.sub.2 ratio, as described in ASTM D-1238. For linear polyolefins, especially linear polyethylene, it is well known that as M.sub.w /M.sub.n increases, I.sub.10 /I.sub.2 also increases.\nJohn Dealy in \"Melt Rheology and Its Role in Plastics Processing\" (Van Nostrand Reinhold, 1990) page 597 discloses that ASTM D-1238 is employed with different loads in order to obtain an estimate of the shear rate dependence of melt viscosity, which is sensitive to weight average molecular weight (M.sub.w) and number average molecular weight (M.sub.n).\nBersted in Journal of Applied Polymer Science Vol. 19, page 2167-2177 (1975) theorized the relationship between molecular weight distribution and steady shear melt viscosity for linear polymer systems. He also showed that the broader MWD material exhibits a higher shear rate or shear stress dependency.\nRamamurthy in Journal of Rheology, 30(2), 337-357 (1986), and Moynihan, Baird and Ramanathan in Journal of Non-Newtonian Fluid Mechanics, 36, 255-263 (1990), both disclose that the onset of sharkskin (i.e., melt fracture) for linear low density polyethylene (LLDPE) occurs at an apparent shear stress of 1-1.4.times.10.sup.6 dyne/cm.sup.2, which was observed to be coincident with the change in slope of the flow curve. Ramamurthy also discloses that the onset of surface melt fracture or of gross melt fracture for high pressure low density polyethylene (HP-LDPE) occurs at an apparent shear stress of about 0.13 MPa (1.3.times.10.sup.6 dynes/cm.sup.2).\nKalika and Denn in Journal of Rheology, 31, 815-834 (1987) confirmed the surface defects or sharkskin phenomena for LLDPE, but the results of their work determined a critical shear stress of 2.3.times.10.sup.6 dyne/cm.sup.2, significantly higher than that found by Ramamurthy and Moynihan et al.\nInternational Patent Application (Publication No. WO 90/03414) published Apr. 5, 1990, discloses linear ethylene interpolymer blends with narrow molecular weight distribution and narrow short chain branching distributions (SCBDs). The melt processibility of the interpolymer blends is controlled by blending different molecular weight interpolymers having different narrow molecular weight distributions and different SCBDs.\nExxon Chemical Company, in the Preprints of Polyolefins VII International Conference, page 45-66, Feb. 24-27 1991, disclose that the narrow molecular weight distribution (NMWD) resins produced by their EXXPOL.TM. technology have higher melt viscosity and lower melt strength than conventional Ziegler resins at the same melt index. In a recent publication, Exxon Chemical Company has also taught that NMWD polymers made using a single site catalyst create the potential for melt fracture (\"New Specialty Linear Polymers (SLP) For Power Cables,\" by Monica Hendewerk and Lawrence Spenadel, presented at IEEE meeting in Dallas, Tex., September, 1991).\nPreviously known narrow molecular weight distribution linear polymers disadvantageously possessed low shear sensitivity or low I.sub.10 /I.sub.2 value, which limits the extrudability of such polymers. Additionally, such polymers possessed low melt elasticity, causing problems in melt fabrication such as film forming processes or blow molding processes (e.g., sustaining a bubble in the blown film process, or sag in the blow molding process etc.). Finally, such resins also experienced melt fracture surface properties at relatively low extrusion rates thereby processing unacceptably."}
{"text": "This invention relates to underwater signaling devices and emergency signaling devices used by scuba divers to perform multiple tasks in a diving emergency. SCUBA is an acronym referring to divers using a Self-Contained Underwater Breathing Apparatus.\nPrior Art Underwater Signaling Devices in General.\nUnderwater signaling devices have been used in scuba diving for years. Some underwater signaling devices use air from the scuba tank which produces an audible low-frequency horn-like sound. Others are designed to bang against the outside of the scuba tank which produces an audible low-frequency clanking-type sound. Still others have been designed to produce a low rattle-type or clicking sound with very limited range.\nA Key Need: Signaling Devices Tailored for Emergency Use.\nThese prior art existing signaling devices are general purpose in nature, and, without a great deal of pre-dive planning, their use in no way indicates an emergency situation in progress under water. Furthermore, they simply are not designed with the goal of effective management of an emergency situation.\nTo understand this, it must first be understood that sensory limitations, specifically visual and auditory limitations, make the scuba diving experience a somewhat isolating one. For example, because of the visual limitations inherent in scuba diving, divers typically pay attention only to their immediate surroundings, in a radius of ten to fifty feet around them, and their safety focus is therefore upon themselves and their diving buddy, who is ordinarily very close by.\nIn addition, auditory limitations, specifically the inability of other divers to hear the human voice, even when a diver is screaming or yelling for help underwater just a few feet away, force divers to rely almost entirely upon their limited sense of sight, specifically line-of-sight, to keep track of their dive buddy, and this is a serious limitation even in good visibility situations.\nStudies of the nature and causes of underwater emergencies lend further support to the critical need for effective communication between and among divers during emergency situations. Statistics indicate that causes of dive emergencies fall in to several categories including: (a) separation from dive buddy (e.g., loss of visual contact due to water clarity; loss of visual contact due to obstructions and distance); (b) equipment failure; (c) underwater entrapment; (d) depletion of air supply; (e) panic; and (f) physical causes (e.g., injuries, illness, fatigue).\nThe situation is further complicated by the following unfortunate reality: notwithstanding the dictates of prudence, many divers simply do not carefully pre-plan audible emergency signal meanings and protocols with other divers in advance, and most carry no signaling device at all. This not only delays recognition of an emergency situation, in addition, it delays coordination of a response by the group. Group coordination is critical, because multiple divers may be necessary to address the emergency under water, and most dive boat captains are precluded from leaving the location of a dive without all divers in hand (and, yes, lives are lost due to delays in retrieving the entire starting contingent of divers).\nFinally, even when an emergency is recognized, current diving emergency practices are frequently limited to the ineffectual use of underwater lights and hand signals as means of visual communication.\nAgainst this backdrop is arrayed the current state of the art in underwater signaling devices, as described above. Tank bangers, air horns, and the like are useful, but these devices are all-purpose signaling devices which are typically used as an ad hoc means of merely attracting the casual, local, non-emergency attention of a single near-by diver. The sound produced by these all-purpose devices has no predetermined meaning to any diver in the water. As a result, other divers in the water regularly ignore these sounds, especially when they are perceived as being far away and not directed specifically to gain that specific diver's attention.\nThe problem, concisely stated, is this: prior art signaling devices are not specifically designed to address the signaling needs unique to underwater emergency situations. Furthermore, even among the various organizations whose missions include the development of improvements to diver safety oriented equipment and procedures, there is no standard signal, signaling protocol, or signaling device that specifically alerts divers to an underwater emergency.\nAfter Ascent: the Need for Accurate Dive Information at the Surface.\nFurthermore, emergencies don't end upon ascent. Frequently, after emergency ascents to the surface, a new problem arises. Divers forced to ascend rapidly due to an emergency may need to undergo recompression treatment, and there is therefore a critical need for comprehensive information about the diver's time under water from the onset of, during, and throughout the resolution of the emergency.\nAlthough there is an increasing trend in diving practices towards the use of personal dive computers and electronic pressure gauges to provide the diver with normal diving information, the limitations of these systems in the face of a need for a clear, easily accessible dive record becomes acutely apparent. Prior art dive computers are simply not designed to record or log, for future use, emergency-related dive information. Also, data retrieval from these systems, in most cases, requires a working knowledge of the unit (and its proprietary user interface) and/or access to the proper hardware, software, and interfaces necessary to extract the data. Thus, valuable time is lost again as divers endeavor to extract critical information from dive computers with which they are not familiar. Finally, these devices typically stop recording data when they are brought to the surface.\nIt is to these underwater emergency management needs, and other circumstances in which expedited, efficient notification is required, that the instant invention is directed."}
{"text": "The present invention relates to a monomeric composition and a polymer obtained by the polymerization thereof or, more particularly, to a monomeric composition capable of being polymerized into a curable polymer which gives a rubbery elastomer having excellent heat and cold resistance and oil resistance as well as a polymer obtained from the monomeric composition.\nSo-called acrylic rubbers belong to a class of synthetic rubbers obtained by the copolymerization of an acrylic monomer such as ethyl acrylate, n-butyl acrylate, 2-methoxyethyl acrylate, 2-ethoxyethyl acrylate and the like with a comonomer which gives crosslinking points in the molecules of the copolymer. Acrylic rubbers obtained from ethyl acrylate as the principal comonomer have excellent oil resistance and heat resistance but the cold resistance thereof is poor. Acrylic rubbers obtained from n-butyl acrylate as the principal comonomer have excellent heat and cold resistance but the oil resistance thereof is poor. Further, acrylic rubbers obtained from 2-methoxyethyl or 2-ethoxyethyl acrylate as the principal comonomer have excellent cold resistance and oil resistance but they have rather poor heat resistance. Thus, none of the conventional acrylic rubbers satisfies the requirements for the oil resistance and cold resistance simultaneously."}
{"text": "Telephony services are offered through a variety of avenues, such as landline phones, cellular phones, and more recently, Voice-Over-Internet Protocol (VoIP). VoIP is a relatively new telephony service that provides communications using Internet protocols rather than the traditional telephone service. VoIP service allows calls to be made to, and received from, traditional phone numbers using a high-speed (broadband) Internet connection (i.e., DSL, cable modem, or broadband wireless technology) instead of using the traditional telephone communication lines. VoIP is implemented by either placing an adapter between a traditional phone and broadband connection or by using a special VoIP phone that connects to a computer or Internet connection.\nTraditional wireline phone services have generally associated a particular phone number with the fixed physical location of the corresponding telephone line. Cellular telephony services determine a cellular caller's physical location by associating the cellular phone with the physical location of the cellular network antenna with which the cellular customer's radio (telephone) is communicating. VoIP services, however, enable consumers to take their home or business phone service almost anywhere because VoIP services can be used from virtually any broadband connection anywhere in the world. This portability raises a number of challenges for the emergency services community in that it makes determining the location of a VoIP caller extremely difficult, if not impossible, because the only information transmitted across the Internet from VoIP callers is the Internet Protocol (IP) addresses associated with the call traffic. Currently there is not a reliable means of mapping an IP address to a precise physical location anywhere in the world.\nEmergency 911 calls from a traditional telephone are usually sent to emergency service providers who are responsible for helping people in a particular geographic area or community. These emergency service providers often can automatically identify the caller's location and direct the closest emergency personnel to that location. They also often automatically identify the caller's telephone number so that the caller can be reached in the event the emergency call is disconnected.\nConsumers with VoIP telephone service have experienced problems accessing 911 emergency services in the same manner as traditional communication services due to the lack of any physical location information associated with their telephone. In some instances, these problems were caused by the consumer failing to provide certain information (such as physical location information) to their VoIP provider in order for their VoIP provider to be able to set up 911 service or the consumer moving their VoIP service to another location without updating their physical location information with the VoIP service provider.\nOne proposed approach for alleviating these problems involves a communication cycle between a VoIP telephone adapter and a communication network where the VoIP telephone adapter communicates with the network once every twenty-four hours. Every time the telephone adapter is disconnected from a power source and reconnected, it communicates with the network to “check in,” and the communication cycle is reset. The network then identifies that the cycle has been reset and recognizes that this could mean that the VoIP telephone adapter, and consequently the customer, has changed locations. Therefore, upon detection of a shift in the communication cycle, the network temporarily suspends the customer's service and posts a message at the customer's web portal directing the customer to confirm the existing registered physical location information or to register a new physical address. Any calls attempted before this physical location information has been confirmed or changed are intercepted and require the caller to confirm or change the physical location information before a call can be completed. This suspension of service does not affect 911 calls, which continue to be associated with the previously registered physical location information.\nThis approach is dependent upon the timing of the communication signal from the telephone adapter and burdens the consumer by denying phone service until the physical location information is resolved. Further, this approach is initiated, for example, every time a customer powers on/off their system, loses power, reboots the router, or loses Internet connectivity; therefore the customer loses phone service after each of these occurrences, which in many cases are not the result of a location change and instead create a “false positive” trigger. This inconveniences the customer by suspending telephone service unnecessarily.\nThe above and other difficulties continue to present challenges to providing effective emergency 911 telephone services and protecting the public safety."}
{"text": "Firearm marksmen, particularly military sharp shooters, have a need for supporting the forward end of a firearm in a stable adjustable manner. Often, a bipod support is used for such front end firearm support. Military sharp shooters have a particular need for a portable, light weight and retractable bipod which also offers significant degrees of adjustability. In particular, it would be useful to have a bipod support having pivotably mounted legs wherein the legs may be adjusted to various positions including a retracted position in which the legs are generally parallel to the longitudinal axis of the firearm. It would also be useful for the legs of such a bipod to have adjustable telescoping portions for adjusting the length of the legs. Moreover, it would be useful if such a bipod support were adapted to allow pivoting adjustment about a vertical axis and a horizontal axis with respect to the legs of the bipod for aiming adjustment."}
{"text": "This invention relates to an improved elevator speed control apparatus for regulating the running speed of an elevator cage to accommodate load changes when passengers exit before the cage comes to a complete halt at an accessed floor.\nA conventional elevator speed control system is shown in FIG. 1, wherein an electric power converter 2 which comprises a plurality of thyristors connected in a 3-phase bridge configuration, is coupled to a 3-phase AC power source 1 and generates DC power that is supplied to an armature 3 of a DC elevator drive motor through a line 2a. The field winding for the motor is not shown in the drawing.\nA tachometer generator 4 driven by the armature 3 produces a speed signal on line 4a proportional to the rotation speed of the armature. A traction sheave 5 also driven by the armature 3 drives an elevator cage 7 and a counterweight 8 through a main cable 6 as is well known. A speed arithmetic circuit 10 receives the speed signal on line 4a from the tachometer generator 4 and a speed instruction signal on line 9a from a speed instruction signal generator 9 as inputs, and generates a current instruction signal on output line 10a. The speed arithmetic circuit 10 along with the speed instruction signal generator 9 and the tachometer 4 constitute a speed control system.\nA current arithmetic circuit 12 receives as inputs the current instruction signal on line 10a from the speed arithmetic circuit 10 and a current signal on line 11a from a current detector 11 proportional to the current supplied to the converter 2. A phase shifter 13 receives the output signal on line 12a from the current arithmetic circuit 12 as an input, and outputs a firing control signal on line 13a for the converter 2. The current arithmetic circuit 12 along with the current detector 11 constitute a current control system.\nBy controlling the firing angle or phase of the thyristor converter 2 by means of both the speed control and current control systems, the voltage applied to the armature 3 is correspondingly controlled and thus the running speed of the elevator cage 7 is controlled through the traction sheave 5. In other words, the elevator cage 7 is speed controlled in accordance with the difference between the speed instruction signal on line 9a and the actual speed signal on line 4a with a high degree of accuracy.\nIn the aforementioned speed control system, in order to compensate for the non-linearity of the converter 2, the response time of the minor loop constituted by the current control system is set at an extremely short value, generally in the range of 0.01 to 0.03 second. On the other hand, the response time of the main loop constituted by the overall speed control system must be set at a higher value in order to avoid resonances in the suspension and traction cables. Consequently, the speed control system is generally designed so as to have a response time in the range of 0.2 to 0.33 second.\nWith such a conventional elevator system, in order to improve the transport efficiency and speed up the overall operation both the internal cage door and the external door on the accessed floor are sometimes controlled to be simultaneously opened just before the cage reaches the floor. A brake system (not shown) is also provided to engage the traction sheave 5, but such engagement does not occur until the cage comes to a complete stop. Passengers may thus step out of the cage before the brake system acts upon the traction sheave, and as a result an abrupt variation in torque is exerted on the sheave due to the change in the cage load or weight, as shown in FIG. 2(a).\nUpon the occurrence of a torque variation, the current flowing through the armature 3 correspondingly varies in response to the output of the speed arithmetic circuit 10 due to the functioning of the current control system, as described above. In this case, however, based on the relatively slow response time characteristics of the speed control system the current flowing through the armature 3 varies or adjusts relatively slowly as shown in FIG. 2 (b). The running speed of the cage 7 therefore varies as shown in FIG. 2 (c), as a result of which the cage may overshoot or undershoot the exact position of the accessed floor, which constitutes a potentially dangerous situation. Even in the best case where the cage ultimately stops at the exact position of the floor sill, the passengers will experience a discomforting \"acceleration-deceleration bump\".\nIt will be understood that the curves of FIG. 2 have been simplified by removing or subtracting therefrom the normal transient values, to leave just the \"abnormal\"variants caused by a premature passenger exit (or entry)."}
{"text": "Glycoprotein hormones, especially those synthesized and secreted by the anterior pituitary gland, play critically important roles in a variety of bodily functions including: metabolism, temperature regulation, growth, and reproduction. The pituitary glycoproteins, luteinizing hormone (LH), follicle stimulating hormone (FSH), and thyroid stimulating hormone (TSH) are similar in structure to the placental gonadotropin, human chorionic gonadotropin (hCG). Each of the molecules is actually a dimer consisting of two protein chains held together by non-covalent, ionic interactions. The alpha chain for each of the hormones is identical. The beta chain is the hormone-specific portion of the dimer.\nFollowing secretion, the hormones travel in the blood stream to the target cells which contain membrane bound receptors. The hormone binds to the receptor and stimulates the cell. Typically, such stimulation involves an increase in activity of a specific intracellular regulatory enzyme which in turn catalyzes a biochemical reaction essential to the response of the cell. In the case of hCG, binding to the hCG receptor present upon the corpus luteum (an ovarian structure), stimulates the activity of the enzyme adenylate cyclase. This enzyme catalyzes the conversion of intracellular ATP to cyclic AMP (cAMP). cAMP stimulates the activity of other enzymes involved in the production of ovarian steroid hormones, especially progesterone. hCG-stimulated progesterone secretion is essential for the maintenance of pregnancy during the first trimester of gestation.\nThe exact mechanism by which a dimeric glycoprotein hormone, such as hCG, stimulates post-receptor events, such as activation of adenylate cyclase activity, is unknown. By a variety of experimental manipulations, it has been shown, however, that the carbohydrate structures, each attached to the hCG molecule by a linkage at respective asparagine residues (N-linked), play important roles in this regard. Treatment of glycoprotein hormones, such as LH, FSH, or hCG with hydrogen fluoride removes approximately 70 percent of the oligosaccharide side chains. The resultant partially deglycosylated molecules retain their receptor binding activity but are unable to stimulate any post-receptor events. Thus it is clear that the sugar portion of the glycoprotein hormone, while not directly involved with receptor binding, plays a critical role in post-receptor events, and therefore, biologic activity.\nIt is also known that in addition to a role in the in vitro bioactivity, oligosaccharides are important components of the molecule's survival time in circulation. Indeed, plasma half-life of a glycoprotein hormone is directly related to the amount of one particular sugar molecule, sialic acid, generally present upon the most distal portion of the oligosaccharide chain. The removal of the carbohydrate portions (by hydrogen fluoride treatment) would result in the production of hormones that bind to the receptor but fail to exert the expected biologic response. In addition, these molecules would have an extremely short plasma half-life since the lack of terminal sialic acid residues would increase the binding affinity to the hepatic asialoglycoprotein receptor thereby hastening clearance from the systemic circulation.\nMany clinical endocrinopathies are the result of over production of stimulating hormones (e.g., excess TSH secretion resulting in hyperthyroidism). A conventional treatment for a pathologic state caused by a hormone excess would be the administration of a hormone antagonist. To be effective, an antagonist must bind with high affinity to the receptor but fail to activate post receptor events. From the earlier discussion, it can be anticipated that hydrogen fluoride treated hormones (that have had approximately 70 percent of carbohydrates removed) would be effective, competitive antagonists. Indeed, it has been shown that HF-treated hormones bind with somewhat greater affinity for the biologic receptor, compete effectively with native material, and diminish the expected biologic action of native hormone in a dose-dependent fashion.\nAt first glance, such a hormone preparation would appear to be a viable candidate for a competitive antagonist therapeutic agent. However, four problems are associated with large scale production of such preparations. First, all pituitary hormone preparations are generally contaminated with other hormones. Thus, while it may be possible to obtain a preparation of a hormone and partially deglycosylate it; the resultant preparation would also contain other deglycosylated hormones as contaminants which may disadvantageously produce unwanted and unacceptable side effects following administration. Secondly, preparations of partially purified hormones vary greatly in their potency, physicochemical characteristics and purity. Therefore, each batch produced would need to be analyzed carefully. The possibility of batch-to-batch variability would necessitate repeated execution of clinical trials to determine the effective dose. Thirdly, the method used to deglycosylate is incomplete and, therefore, somewhat uncontrollable. No information is available as to the variability associated with hydrogen fluoride treatment of successive batches of hormone. Potential variability in this process would also require extensive characterization of each batch produced. Fourth, carbohydrate side chains also play an important role in dictating the plasma half-life of a molecule. Partially deglycosylated hormones have been shown to be rapidly cleared from the circulation following injection. Therefore, repeated injections of large doses of deglycosylated hormones would be required to achieve a desired effect. The need for such large doses of hormone to deglycosylate and administer creates yet another problem, namely availability. Large quantities of hormones are presently unavailable and could conceivably only be made available through recombinant DNA technology.\nAs mentioned above, deglycosylated hormones, while exhibiting the desired competitive antagonistic properties in vitro would be of little therapeutic value in vivo due to their extremely short plasma half-life. The only potential solution to this dilemma would be to preferentially remove carbohydrate residues that are responsible for imparting the molecules' biologic action and sparing those that provided the molecules' long plasma half-life. It is, however, impossible to obtain this result with conventional chemical or enzymatic means (i.e. hydrogen fluoride treatment or enzymatic digestion) because of the non-specific nature of the chemical treatment.\nIt is an object of the present invention to provide a method for preferentially removing carbohydrate residues that are responsible for imparting biologic action to molecules while sparing those associated with a long plasma half-life.\nIt is yet another object of the present invention to provide a method for obtaining molecules via non-chemical treatment which have the desired antagonistic nature coupled with a long half-life.\nIt is still another object of the present invention to provide a recombinant technique for obtaining therapeutically effective glycoprotein hormones having been partially deglycosylated by having at least one oligosaccharide chain entirely removed.\nIt is still yet another object of the present invention to provide novel hormonal competitive antagonists having therapeutic utility.\nIt is yet still another object of the present invention to provide competitive hormone antagonists having substantially batch-to-batch uniformity and consistent potency."}
{"text": "Sutures are available in a variety of materials, shapes and sizes. One such shape is a “flat” suture, where a plurality of fibers are braided such that the resulting suture is significantly wider than it is tall, resulting in a flat, or planar shape. Such sutures can be useful, for example, in increasing the contact surface area between the suture and the underlying soft tissue to help increase the contact footprint of the soft tissue against the underlying bone. This is particularly useful in, for example, rotator cuff repairs where a “suture bridge” is formed over the upper surface of the cuff to compress the cuff tissue to the underlying bone. Another advantage of such flat sutures is that the larger surface area distributes forces exerted on the tissue by the suture such that there is less of a chance the suture will cut into the tissue relative to a thinner suture (e.g., a traditional round suture).\nCurrent flat sutures on the market, however, suffer from multiple drawbacks such that surgeons have been slow to utilize them. For example, the shape of the flat suture often times requires different instrumentation, thereby rendering much of the instrumentation on the market incompatible with such flat sutures. Commonly, surgeons have been using a particular instrument for a long time and are unwilling to change just to accommodate a differently-shaped suture. Another drawback of current flat sutures occurs when a surgeon attempts to tie the flat suture in a knot. Typically, the flat shape of the braided suture does not compress well, and thus the resulting knot easily loosens, or worse, comes completely undone.\nThus, there is a need in the art for a flat suture that is compatible with instrumentation on the market and which can be manipulated and handled effectively by traditional methods, such as when forming a secure knot."}
{"text": "The present, invention relates to a manufacturing method of a combination material of metal foil and ceramic by joining a metal foil onto surfaces of various types of ceramics, and also relates to a metal foil laminated ceramic substrate manufactured from said combination material of metal foil and ceramic.\nRecently, a new improvement has been attempted by combining ceramics with metals in the application fields of ceramics. For example, a metal material composed of nickel alloy, titanium alloy, chromium alloy, or the like is joined to a ceramic material composed of alumina, zirconia, magnesia, or the like by using a combination technique such as diffusion combination so as to manufacture a combination material of metal and ceramic, which is used in various devices and apparatuses.\nHowever, since such conventional combination material of metal and ceramic is manufactured by simply joining a metal material to a ceramic material under a certain pressure, when the pressure is applied, there sometimes occurs fracture of the ceramic material which is a much more brittle material as compared with the metal material, thus causing problems in productivity.\nAn object of the present invention is to solve the problem mentioned above and provide a manufacturing method of a combination material of metal foil and ceramic capable of completely joining a metal foil to a ceramic material even with a low pressure applied thereto and preventing the ceramic material from fracturing so as to improve the productivity of a combination material of metal foil and ceramic or a metal foil laminated ceramic substrate, and also provide such metal foil laminated ceramic substrate.\nIn order to achieve the object mentioned above, the manufacturing method of a combination material of metal foil and ceramic comprises the steps of ion-etching a surface of a metal foil and a ceramic material to be joined together to activate and clean the surfaces, heating said surface of the ceramic material, which is held on a holder, to a temperature range of 250 to 500xc2x0 C., pressure-welding said surface of the metal foil to said surface of the ceramic material held on the holder under a pressure not more than 1 kg/mm2, and thus heat-joining the metal foil and the ceramic material to manufacture a combination material of metal foil and ceramic."}
{"text": "Instrument clusters on automobiles generally include a plurality of gauges for displaying such operational information such as vehicle speed, engine RPM, engine temperature, fuel level and many other information. The gauges may include analog or digital readings for displaying the information depending on manufacturer and styling preferences. An analog gauge typically includes a faceplate having indicia thereon such as numbers and a pointer for rotating to the appropriate number.\nOne important design consideration for an instrument cluster and related gauges is the ability of a vehicle operator to easily view and read the gauges in all driving environments. In particular, nighttime driving requires the instrument cluster to illuminate in some fashion whereby the numbers and corresponding pointers are easily distinguishable."}
{"text": "1. Technical Field\nThe present disclosure relates generally to medical catheters, and more particularly to a urology catheter including an advantageous tip configuration and related method of manufacture.\n2. Description of the Related Art\nAs is known, urinary catheters may be employed to transport urine collected in the bladder out of a patient via the urinary tract. For example, urinary catheters such as Foley catheters have a shaft including a drainage lumen that communicates with drainage eyes disposed adjacent a distal end thereof. An inflatable balloon is disposed adjacent the distal end of the shaft. During placement, the distal end of the shaft is passed through the patient's urethra until the balloon and drainage eyes are located in the patient's bladder. The balloon is inflated through an inflation lumen to retain the catheter in the bladder. Urine may drain through the drainage eyes and drainage lumen, which is in communication with a proximal end of the catheter. The proximal end of the catheter is connected to a receptacle for collection of urine.\nVarious known urology catheters may be produced with a liquid injection molded tip that is attached to the catheter shaft. Typically, the drain eyes of the urinary catheter are formed via a punching operation. This punching operation, however, can disadvantageously result in jagged edges around the eyes and an inconsistent eye size. Other known urinary catheters may include molded drain eyes.\nThese known manufacturing apparatus and techniques suffer from various drawbacks. For example, the size of the drain eye may be limited disadvantageously resulting in inferior drainage or rapid blockage due to encrustation. Alternatively, due to manufacturing constraints, the size of the drain eye may have to be increased to accommodate the manufacturing apparatus. An increase in size of the drain eye may cause the tip to become undesirably flexible, disadvantageously resulting in difficult insertion with the body due to tip deflection.\nTherefore, it would be desirable to have a urology catheter including an advantageous tip configuration and related method of manufacture that improves drainage characteristics, avoids encrustation and enhances stiffness to facilitate tip insertion with a patient. Desirably, a drain eye of the urology catheter is defined by an external radius and an internal radius of the tip. It would be highly desirable if the urology catheter and its constituent parts are easily and efficiently manufactured and assembled."}
{"text": "This invention generally relates to a garden tool, and more particularly to an improved tool for removing individual weeds and unwanted plants from the earth by directly engaging their roots without separating the roots from the stem and with minimal disturbance of surrounding earth and other plants."}
{"text": "Biomass thermal conversion is an attractive method for generating synthetic gas to run engines or to produce useful end products such as charcoal. Carbonaceous byproducts are typically inexpensive or free to source. Unfortunately, biomass byproducts come in a wide array of shapes and sizes, and extra machinery is usually required to pre-process the feedstock into forms acceptable to gasification or pyrolysis machines. This processing equipment is often expensive and difficult to operate, which challenges the ultimate attractiveness of biomass thermal conversion solutions.\nThus, there is a need in the field of biomass thermal conversion for system capable of utilizing a wide range of fuel shapes and sizes, without feedstock preprocessing on the front end. This invention provides such a solution through a novel “reactor-internal” fuel processing solution that reduces a wide range of input biomass feedstock to a common size of granulated char."}
{"text": "The present invention relates to a microfluidic device, which can be interfaced to a mass spectrometer (MS). The device comprises a microchannel structure having a first port (inlet port) and a second port (outlet port). A sample to be analysed is applied to the first port and presented to the mass spectrometer in the second port. This second port will be called an MS-port. There may be additional inlet and outlet ports. During passage through the microchannel structure the sample is prepared to make it suitable for analysis by mass spectrometry.\nThe sample presented in an MS-port will be called an MS-sample. An analyte in an MS-sample is an MS-analyte. xe2x80x9cSamplexe2x80x9d and xe2x80x9canalytexe2x80x9d without prefix will primarily refer to a sample applied to an inlet port.\nOne important aspect of the present invention concerns mass spectrometry in which the MS-samples are subjected to Energy Desorption/Ionisation from a surface by input of energy. Generically this kind of process will be called EDI and the surface an EDI surface in the context of the invention. Typicallly EDIs are thermal desorption/ionisation (TDI), plasma desorption/ionisation (PDI) and various kinds of irradiation desorption/ionisation (IDI) such as by fast atom bombardment (FAB), electron impact etc. In the case a laser is used the principle is called laser desorption/ionisation (LDI). Desorption may be assisted by presenting the MS analyte together with various helper substances or functional groups on the surface. Common names are matrix assisted laser desorption/ionisation (MALDI) including surface-enhanced laser desorption/ionisation (SELDI). For MALDI see the publications discussed under Background Publications below. For SELDI see WO 0067293 (Ciphergen Biosystems).\nThe surface from which desorption/ionisation is intended to take place is called an EDI surface.\nBy microformat is meant that in least a part of the microchannel structures the depth and/or width is in the microformat range, i.e.  less than 103 xcexcm, preferably  less than 102 xcexcm. In the most typical microformat structures either the width and/or the depth are in principle within these ranges essentially everywhere between the sample inlet port and the MS-port.\nFor some time there has been a demand for microfluidic sample handling and preparation devices with integrated MS-ports. This kind of devices would facilitate automation and parallel experiments, reduce loss of analyte, increase reproducility and speed etc.\nWO 9704297 (Karger et al) describes a microfluidic device that has an outlet port that is claimed useful when conducting electrospray ionisation mass spectrometry (ESI MS), atmospheric pressure chemical ionisation mass spectrometry (APCI MS), matrix assisted laser desorption/ionisation mass spectrometry (MALDI MS) and a number of other analytical principles.\nU.S. Pat. No. 6,110,343 (Ramsey et al) describes an electrospray interface between a microfluidic device and a mass spectrometer.\nU.S. Pat. No. 5,969,353 (Hsieh) describes an improved interface for electrospray ionization mass spectrometry. The interface is in the form of an electrospray tip connected to a microchannel structure of a chip.\nU.S. Pat. No. 5,197,185 (Yeung et al) describes a laser-induced vaporisation and ionization interface for directly coupling a microscale liquid based separation process to a mass spectrometer. A light-adsorbing component may be included in the eluting liquid in order to facilitate vaporisation.\nU.S. Pat. No. 5,705,813 (Apffel et al) and U.S. Pat. No. 5,716,825 (Hancock et al) describe a microfluidic chip containing an interface between a microfluidic device and an MALDI-TOF MS apparatus. The microfluidic device comprises\n(a) an open ionisation surface that may be used as the probe surface in the vaccum gate of an MALDI-TOF MS apparatus (column 6, lines 53-58 of U.S. Pat. No. 5,705,813) or\n(b) a pure capture/reaction surface from which the MS-analyte can be transferred to a proper probe surface for MALDI-TOF MS (column 12, lines 13-34, of U.S. Pat. No. 5,716,825).\nThese publications suggest that means, such as electrical connections, pumps etc, for transporting the liquid within a microchannel structure of the device are integrated with or connected to the device. This kind of transporting means imposes an extra complexity on the design and use, which in turn may negatively influence the production costs, easiness of handling etc of these devices.\nU.S. Pat. No. 5,705,813 (Apffel et al) and U.S. Pat. No. 5,716,825 (Hancock et al) are scarce about\nthe proper fluidics around the MALDI ionisation surface,\nthe proper crystallisation on the MALDI ionisation surface,\nthe proper geometry of the port in relation to crystallisation, evaporation, the incident laser beam etc,\nthe proper arrangement of conductive connections to the MALDI ionisation surface for MALDI MS analysis.\nWO 04297 (Karger et al) and WO 0247913 (Gyros AB) suggest to have microchannel structures in radial or spoke arrangement.\nA number of publications referring to the use of centrifugal force for moving liquids within microfluidic systems have appeared during the last years. See for instance WO 9721090 (Gamera Bioscience), WO 9807019 (Gamera Bioscience) WO 9853311 (Gamera Bioscience), WO 9955827 (Gyros AD), WO 9958245 (Gyros AD), WO 0025921 (Gyros AB), WO 0040750 (Gyros AB), WO 0056808 (Gyros AB), WO 0062042 (Gyros AD) and WO 0102737 (Gyros AB) as well as WO 0147637 (Gyros AB), WO 0154810 (Gyros AB), WO 0147638 (Gyros AB), and WO 01411465.\nSee also Zhang et al, xe2x80x9cMicrofabricated devices for capillary electrophoresisxe2x80x94electrospray mass spectrometryxe2x80x9d, Anal. Chem. 71 (1999) 3258-3264) and references cited therein.\nKido et al., (xe2x80x9cDisc-based immunoassay microarraysxe2x80x9d, Anal. Chim. Acta 411 (2000) 1-11) has described microspot immunoassays on a compact disc (CD). The authors suggest that a CD could be used as a continuous sample collector for microbore HPLC and subsequent detection for instance by MALDI MS. In a preliminary experiment a piece of a CD manufactured in polycarbonate was covered with gold and spotted with a mixture of peptides and MALDI matrix.\nA first object is to provide improved means and methods for transporting samples, analytes including fragments and derivatives, reagents etc in microfluidic devices that are capable of being interfaced with a mass spectrometer.\nA second object is to provide improved microfluidic methods and means for sample handling before presentation of a sample analyte as an MS-analyte. Sub-objects are to provide an efficient concentration, purification and/or transformation of a sample within the microfluidic device while maintaining a reproducible yield/recovery, and/or minimal loss of precious material.\nA third object is to provide improved microfluidic methods and means that will enable efficient and improved presentation of the MS-sample/MS-analyte. This object in particular applies to MS-samples that are presented on a surface, i.e. an EDI surface.\nA fourth object is to enable reproducible mass values from an MS-sample that is presented on a surface, i.e. on an EDI surface.\nA fifth object is to provide improved microfluidic means and methods for parallel sample treatment before presentation of the analyte to mass spectrometry. The improvements-of this object refer to features such as accuracy in concentrating, in chemical transformation, in required time for individual steps and for the total treatment protocol etc. By parallel sample treatment is meant that two or more sample treatments are run in parallel, for instance more than five, such as more than 10, 50, 80, 100, 200, 300 or 400 runs. Particular important numbers of parallel samples are below or equal to the standard number of wells in microtiter plates, e.g. 96 or less, 384 or less, 1536 or less, etc\nA sixth object is to provide a cheap and disposable microfluidic device unit enabling parallel sample treatments and having one or more MS-ports that are adapted to a mass spectrometer.\nThe present inventors have recognized that several of the above-mentioned objects can be met in the case inertia force is used for transportation of a liquid within a microfluidic device of the kind discussed above. This is applicable to any liquid that is used in the microfluidic device, for instance washing liquids and liquids containing at least one of (a) the analyte including derivatives and fragments thereof, (b) a reagent used in the transformation of the sample/analyte, etc.\nThe present inventors have also recognized that one way of optimizing an EDI area within a microfluidic device is related to\n(a) the design and/or positioning of a conducting layer in the EDI area, and/or\n(b) the importance of a conductive connection to the EDI area for MS analysis.\nThis kind of connection supports the proper voltage and/or charge transport at the EDI area, for instance.\nImproper conductive properties may interfere with the mass accuracy, sensitivity, resolution etc.\nConductive and non-conductive properties shell refer to the property of conducting electricity.\nA first aspect of the invention is thus a method for transforming a liquid sample containing an analyte to an MS-sample containing an MS-analyte and presenting the MS-sample to a mass spectrometer. The method is characterized in comprising the steps of:\n(a) applying the liquid sample to an inlet port of a covered microchannel structure of a microfluidic device,\n(b) transforming the liquid sample to an MS-sample containing the MS-analyte within the microchannel structure, and\n(c) presenting the MS-analyte to the mass spectrometer.\nA further characteristic feature of this aspect is that transport of liquid within the microchannel structure is performed by the application of inertia force. Inertia force may be the driving force in only a part of the microchannel structure or the whole way from an inlet port to an MS-port and/or to any other outlet port. It is believed that the most general and significant advantages of using inertia force will be accomplished in so called transporting zones, i.e. between zones having predetermined functionalities, or for overcoming or passing through valve functions within a microchannel structure (capillary junctions, hydrophobic breaks etc). See below. The MS-port typically has a conductive connection for MS analysis.\nAt the priority date the most important inertia force for microfluidic devices is centrifugal force. In other words a force that causes outward radial transportation of liquid by spinning a disc in which the liquid is located within microchannel structures that are oriented radially (spinning is around an axis that is perpendicular to the plane of the disc). Inertia force caused by other changes of direction and/or magnitude of a force can be utilized.\nThe first aspect also includes the corresponding mass spectrometric method, i.e. the same method together with the actual collection of a mass spectrum and analysis thereof, for instance in order to gain molecular weight and structure information about the analyte.\nThe first aspect is further defined as discussed below for the microfluidic device as such and for the individual steps.\nA second aspect of the invention is a microfluidic device containing one, two or more microchannel structures containing an inlet port, an MS-port and a flow path connected to one or both of the ports. The device may be disc-formed or otherwise provide a planar form. The characteristic feature is that the microchannel structures are oriented radially in an annular/circular arrangement. Thus each microchannel structure extends in a radial direction with an inlet port at an inner position and an outlet port such as an MS-port, at an outer peripheral position. The MS-port typically has a conductive connection as discussed above. The features discussed below further define this aspect of the invention.\nA third aspect of the invention is a microfluidic device comprising a plurality of covered microchannel structures as defined herein and with each microchannel structure having an MS-port comprising an EDI area in which there is a conducting layer (layer I). This aspect of the present invention comprises a number of subaspects having the common characteristic feature that there may be a conductive connection to layer (I) of each individual EDI area, as discussed above. There are also features that are distinct for each subaspect.\nA first subaspect is further characterized in that layer (I) of each EDI area is part of a continuous conducting layer that is common for two or more up to all of the EDI-areas.\nA second subaspect is further characterized in that in each EDI area there is a non-conducting layer (layer II) between layer (I) and the surface of the EDI area. Layer (II) in each EDI area may be part of a continuous non-conducting layer that is common for two or more up to all of the EDI-areas.\nA third subaspect is further characterised in that each MS-port has an opening that is restricted by a lid which is common for and covers a number of microchannel structures. The lid may have a conducting layer that at least embraces the openings that are present in the lid. The conducting layer may be continuous in the sense that it covers at least the areas around and between the openings of two or more up to all of the MS-ports. This layer may have a conductive connection as discussed above.\nA fourth subaspect is similar to the third subaspect in the sense that there is a lid covering at least a part of each microchannel structures. In this subaspect the lid also covers or restricts the openings of the MS-ports and is removable to an extent that enables exposure of the opening in each MS-port, for instance exposing the surfaces of EDI areas. For EDI ports the removal will facilitate irradiation and the desorption/ionisation of the MS-analyte. The removal may also facilitate evaporation of volatile components.\nThe sample applied to an inlet port may contain one or more analytes, which may comprise lipid, carbohydrate, nucleic acid and/or peptide structure or any other inorganic or organic structure. The sample treatment protocol to take place within the microchannel structure typically means that the sample is transformed to one or more MS-samples in which\n(a) the MS-analyte is a derivative of the starting analyte and/or\n(b) the amount(s) of non-analyte species have been changed compared to the starting sample, and/or\n(c) the relative occurrence of different MS-analytes in a sample is changed compared to the starting sample, and/or\n(d) the concentration of an MS-analyte is changed relative the corresponding starting analyte in the starting sample, and/or\n(e) sample constituents, such as solvents, have been changed and/or the analyte has been changed from a dissolved form to a solid form, for instance in a co-crystallized form.\nItem (a) includes digestion into fragments of various sizes and/or chemical derivatization of an analyte. Digestion may be purely chemical or enzymatic. Derivatization includes so-called mass tagging of either the starting analyte or of a fragment or other derivative formed during a sample treatment protocol, which takes place in the microchannel structure. Items (b) and/or (c) include that the sample analyte has been purified and/or concentrated. Items (a)-(d), in particular, apply to analytes that are biopolymers comprising carbohydrate, nucleic acid and/or peptide structure.\nThe sample is typically in liquid form and may be aqueous.\nThe sample may also pass through the microchannel structure without being changed. In this case the structure only provide a proper form for dosing of the analyte to the mass spectrometer."}
{"text": "1. Field of the Invention\nThe present invention relates to an arrangement on a vehicle with a rotatable crane, comprising a crane boom and associated hoist winch and rescue winch. Such vehicles are made in particular as rescue vehicles for civilian and military use, intended for rescuing vehicles--everything from lightweight automobiles to heavy tanks--which have driven off the road or become stuck in difficult terrain.\n2. Description of the Prior Art\nPreviously known rescue vehicles are built on heavy truck chassis with drive on four or more wheels, and they are provided with cranes and winches for hoisting and hauling, and as a rule they also have support feet to improve stability and to prevent the vehicle from sinking down into the ground or surface on which it stands when said surface has less supporting capacity than, e.g., a roadway.\nOn previously known rescue vehicles, a rescue winch is mounted on the vehicle with its axis of revolution ordinarily transverse of the longitudinal axis of the vehicle. This means that the cable can be pulled straight out from the winch along the length of the vehicle without having to pass directional pulleys, guide rolls or the like. In such conditions the fixedly mounted winch will work satisfactorily. If, however, it is necessary to haul in a tank, for example, in a direction other than lengthwise along the rescue vehicle, the cable from the winch must be guided in several rather sharp bends in order to come into the desired direction for pulling. Such sharp bends produce great wear on the cable and increase friction so that the pulling power of the winch cannot be utilized fully. Furthermore, situations may also arise in which one wishes to run the cable from the rescue winch over the crane arm or boom of the rescue vehicle, and in such situations it will be necessary, assuming it is possible at all, to apply further sharp bends on the cable."}
{"text": "Companies and other organizations rely on data feeds from numerous sources to obtain relevant, timely data for their businesses and operations. Such data sources, however, are commonly received in a variety of formats which requires the receiving organization to manually extract the relevant data from each type of document. Consequently, many organizations expend considerable resources to analyze incoming data and file formats, identify relevant data, and store it in a useful format in their existing databases. The extent of effort involved may consume a large amount of the organization's administrative resources. In addition, many organizations do not effectively correlate incoming data with their existing databases. Consequently, organizations frequently miss opportunities to effectively and rapidly assimilate incoming data into their operations and businesses. These and other drawbacks exist with known systems."}
{"text": "1. Technical Field\nThe present invention relates to data processing and, in particular, to scripts in a network data processing system. Still more particularly, the present invention provides a method, apparatus, and program for evaluating scripts.\n2. Description of Related Art\nThe worldwide network of computers commonly known as the “Internet” has seen explosive growth in the last several years. Mainly, this growth has been fueled by the introduction and widespread use of so-called “web browsers,” which enable simple graphical user interface-based access to network servers, which support documents formatted as so-called “web pages.” These web pages are versatile and customized by authors. For example, web pages may mix text and graphic images. A web page also may include fonts of varying sizes.\nA browser is a program that is executed on a graphical user interface (GUI). The browser allows a user to seamlessly load documents from the Internet and display them by means of the GUI. These documents are commonly formatted using markup language protocols, such as hypertext markup language (HTML). Portions of text and images within a document are delimited by indicators, which affect the format for display. In HTML documents, the indicators are referred to as tags. Tags may include links, also referred to as “hyperlinks,” to other pages. The browser gives some means of viewing the contents of web pages (or nodes) and of navigating from one web page to another in response to selection of the links.\nBrowsers may also read and interpret pages including scripts, such as JAVAScript or JScript. JAVAScript is a popular scripting language that is widely supported in Web browsers and other Web tools. JAVAScript adds interactive functions to HTML pages, which are otherwise static, since HTML is a display language, not a programming language. JScript is similar to JAVAScript, but has extensions specifically for the Microsoft Windows environment.\nDifferent Web browser software applications support scripts to different degrees. In fact, different versions of a browser application may support scripts differently. The rivalry between browsers in the market, such as Netscape Navigator by Netscape Communications and Internet Explorer by Microsoft Corporation, has led to a disparity between standards. Netscape uses different syntax with JAVAScript than Internet Explorer uses with JScript.\nBrowser and platform dependent script functions can cause a nightmare for script development, testing, and maintenance. Therefore, it would be advantageous to provide an improved script evaluator for determining browser support."}
{"text": "Titanium-aluminium (Ti—Al) alloys and alloys based on titanium-aluminium (Ti—Al) inter-metallic compounds are very valuable materials. However, they can be difficult and expensive to prepare, particularly in the powder form. This expense of preparation limits wide use of these materials, even though they have highly desirable properties for use in aerospace, automotive and other industries.\nReactors and methods for forming titanium-aluminium based alloys have been disclosed. For example, WO 2007/109847 discloses a stepwise method for the production of titanium-aluminium compounds and titanium alloys and titanium-aluminium inter-metallic compounds and alloys.\nWO 2007/109847 describes the production of titanium-aluminium based alloys via a two stage reduction process, based on reduction of titanium tetrachloride with aluminium. In stage 1, TiCl4 is reduced with Al in the presence of AlCl3 to produce titanium subchlorides according to the following reaction:TiCl4+(1.333+x)Al→TiCl3+(1+x)Al+0.333AlCl3 or  (1)TiCl4+(1.333+x)Al→TiCl2+(0.666+x)Al+0.666AlCl3  (1)\nIn stage 2, the products from reaction (1) are processed at temperatures between 200° C. and 1300° C. to produce a powder of titanium-aluminium based alloys, according to the following (simplified) reaction scheme:TiCl3+(1+x)Al→Ti—Al+AlCl3 or  (2)TiCl2+(0.666+x)Al→Ti—Alx+0.666AlCl3  (2)"}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nFuel cells are useful as a power source for electric vehicles and other applications. An exemplary fuel cell has a membrane electrode assembly (MEA) with catalytic electrodes and a proton exchange membrane (PEM) formed between the electrodes. Water is generated at the cathode electrode based on the electrochemical reactions between hydrogen and oxygen occurring within the MEA. Gas diffusion media plays an important role in PEM fuel cells. Generally disposed between catalytic electrodes and flow field channels that introduce reactant gases into the fuel cell, the gas diffusion media provide pathways for reactant to diffuse to the electrode and pathways for removal of the product water, electronic conductivity, and heat conductivity, as well as mechanical strength needed for proper functioning of the fuel cell.\nDuring operation of the fuel cell, water is generated at the cathode based on electrochemical reactions involving hydrogen and oxygen occurring within the MEA. Efficient operation of a fuel cell depends on the ability to provide effective water management in the system. For example, the diffusion media prevent the electrodes from flooding (i.e., filling with water and severely restricting O2 access) by removing product water away from the catalyst layer while maintaining reactant gas flow from gas flow channels of the bipolar plate through to the catalyst layer.\nFuel cell stacks can contain a large number of fuel cells depending on the power requirement of the application. For example, typical fuel stacks have up to 400 individual fuel cells and more. Because the fuel cells in the stacks operate in series, a weakness or poor performance in one cell can translate into poor performance of the entire stack. For this reason, it is desirable for every fuel cell in the stack to operate at high efficiency.\nTypical manufacturing steps for gas diffusion media include manufacturing a carbon fiber paper, impregnating the paper with resin or a mixture of resin and fillers, molding the impregnated paper, and carbonization or graphitization of the resin-impregnated carbon fiber paper. The steps of manufacturing paper and impregnation are continuous, while the molding, carbonization and graphitization steps may be either batch or continuous.\nBecause typical fuel stacks contain so many individual fuel cells, it is important for the manufacturing process of the diffusion media to have a high degree of reliability. Improvements in the manufacturing process that reduce cost, simplify the process, or enhance performance of the media are thus desirable."}
{"text": "Many implanted medical devices that are powered by electrical energy have been developed. Most of these devices comprise a power source, one or more conductors, and a load.\nWhen a patient with one of these implanted devices is subjected to high intensity magnetic fields, currents are often induced in the implanted conductors. The large current flows so induced often create substantial amounts of heat. Because living organisms can generally only survive within a relatively narrow range of temperatures, these large current flows are dangerous.\nFurthermore, implantable devices, such as implantable pulse generators (IPGs) and cardioverter/defibrillator/pacemaker (CDPs), are sensitive to a variety of forms of electromagnetic interference (EMI). These devices include sensing and logic systems that respond to low-level signals from the heart. Because the sensing systems and conductive elements of these implantable devices are responsive to changes in local electromagnetic fields, they are vulnerable to external sources of severe electromagnetic noise, and in particular to electromagnetic fields emitted during magnetic resonance imaging (MRI) procedures. Therefore, patients with implantable devices are generally advised not to undergo magnetic resonance imaging (MRI) procedures, which often generate static magnetic fields of from between about 0.5 to about 10 Teslas and corresponding time-varying magnetic fields of about 20 megahertz to about 430 megahertz, as dictated by the Lamor frequency (see, e.g., page 1007 of Joseph D. Bronzino's “The Biomedical Engineering Handbook,” CRC Press, Hartford, Conn., 1995). Typically, the strength of the magnetic component of such a time-varying magnetic field is about 1 to about 1,000 microTesla.\nOne additional problem with implanted conductors is that, when they are conducting electricity and are simultaneously subjected to large magnetic fields, a Lorentz force is created which often causes the conductor to move. This movement may damage body tissue.\nIn U.S. Pat. No. 4,180,600, there is disclosed and claimed a fine magnetically shielded conductor wire consisting of a conductive copper core and a magnetically soft alloy metallic sheath metallurgically secured to the conductive core, wherein the sheath consists essentially of from 2 to 5 weight percent of molybdenum, from about 15 to about 23 weight percent of iron, and from about 75 to about 85 weight percent of nickel. Although the device of this patent does provide magnetic shielding, it still creates heat when it interacts with strong magnetic fields.\nIt is an object of this invention to provide a sheath assembly, which is shielded from magnetic fields."}
{"text": "1. Field of the Invention\nThis invention pertains to the general field of optical filters and, in particular, to temperature-stable and tunable high-performance etalon filters.\n2. Description of the Prior Art\nEtalons are well known optical devices that consist of two reflective surfaces parallel to one another and spaced apart by solid spacers to produce a predetermined optical length (the “cavity length”). They may consists simply of a solid parallel plate with reflective surfaces (so called “solid etalons”) or of two plates with an air gap between them that defines the cavity (so called “air-spaced etalons”), as illustrated in FIG. 1. A hybrid form of etalon (so called “re-entrant etalon”) utilizes an additional solid structure with a reflective surface (a “riser”) between the two plates in order to achieve narrower cavity lengths than practically obtainable with the use of spacers.\nWhen illuminated with a broadband collimated light, etalons produce a transmission beam and a reflection beam with periodic spectra characterized by very narrowband spikes centered at wavelength determined by the physical properties and dimensions of the etalon. A typical etalon transmission spectrum is illustrated in FIG. 2. With reference to air-spaced etalons, in particular, the specific center wavelength λ′ of the passband (the spectral spike) and the period between spectral spikes (commonly referred to in the art as channel spacing or free spectral range—FSR—of the device) are a function of the optical length of the etalon's cavity. This disclosure is limited to a discussion of transmission operation because those skilled in the art would readily understand that it is similarly applicable to reflection operation.\nIn particular, referring for example to the etalon 10 and the intensity spectrum 12 of FIGS. 1 and 2, respectively, minor changes in the optical length L (corresponding to the geometric length L′ shown in the figures) of the cavity 14 will cause a shift of the periodic spectrum along the wavelength axis, as indicated by arrows 16. As is well understood by those skilled in the art, varying the optical length L of the cavity also produces a change in the width of the spectral spike and in the free spectral range of the etalon.\nThese properties of etalons are very advantageous for many optical applications. In particular, etalons are used as high-performance filters to isolate light of a very a precise frequency, as may be needed for a particular application. In telescopic astronomy, for instance, such filters are particularly useful for observing objects at specific wavelengths. Since the exact wavelength of each peak is a function of the exact optical length L of the cavity, it has been most important in the art to build etalon filters with precise and uniform spacing between the two plates (18,20) constituting the etalon (FIG. 1). To that end, very precisely machined spacers 22,24 of equal thickness L′ are used, typically uniformly distributed around the annular periphery of the plates in a sufficient number to separate the plates and produce a cavity of uniform optical length L. Moreover, these spacers are typically made of materials having a low coefficient of thermal expansion. (It is noted that L′ is the physical cavity length corresponding to the desired optical path length L, the two quantities being related by the equation L=nL′, where n is the index of refraction of the medium in the cavity.).\nIn practice it has been difficult and expensive to achieve the desired degree of perfection because of the very narrow tolerances (in the order of nanometers) required for the level of performance associated with extremely narrowband applications. U.S. Pat. Nos. 6,181,726 and 6,215,802 disclosed several advances over the prior art whereby the uniformity of the etalon's optical length was improved. According to one approach described in the patents, all the spacers used to form the etalon are selected from a common local area of a spacer substrate produced by standard-precision optical manufacturing techniques. It was discovered that, as a result of this selection, the spacers tend to have substantially more uniform thickness and, therefore, they produce a more uniform etalon cavity. According to another, complementary approach, an additional spacer from the same local substrate area is used at the center of the etalon, thereby providing a correction to plane deformations produced by the optical contact of the peripheral spacers with the etalon plates.\nWhile the techniques described in these patents provide a significant improvement over the etalons previously known in the art, they are very labor-intensive and therefore expensive to practice. In addition, the resulting etalons, while more uniform in the optical length of the cavity, are not necessarily tuned to the precise desired wavelength. Copending U.S. Ser. No. 10/795,167 discloses a solution to this problem based on the use of counterbalanced forces applied to the etalon elements. This approach constitutes another significant advance in the art, but it does not address the problem of changes in performance (in terms of center wavelength and FSR) produced by thermal variations. Therefore, there is still a need for an etalon structure that is relatively insensitive to thermal effects and that produces extremely accurate tuning of the optical length of the etalon cavity. The present invention provides a solution to this remaining challenge that also produces a greater range of angular acceptance and a mechanism for thermally tuning the etalon to a precise level of performance."}
{"text": "In general, this invention relates to photographic film cassettes and, more particularly, an improved pressure pad spring and cassette arrangement therefor which avoids diminishing spring effectiveness and improves overall cassette assembly.\nMultipurpose film cassettes have been developed in which a strip of photographic film is operated so as to be exposed, processed and projected without leaving the cassette. Film cassettes of this type are disclosed in several U.S. patents assigned in common with the present invention.\nIn film cassettes of this category, a supply of light sensitive photographic film can be selectively exposed in a camera particularly adapted to receive and operate the cassette. To process or develop the exposed film, the cassette is removed from the camera and placed in a player or processing and viewing apparatus capable of activating a cassette contained processor for depositing a desired uniform layer of processing fluid on the film's exposed emulsion surface. During such processing a conventional series of successive, positive transparent images on the exposed film is developed. Following processing in the manner indicated, the player apparatus is operated as a projector. During projection, the film is incrementally advanced, frame-by-frame, past a light source. Accordingly, the series of positive transparent images of the scenes to which the film were exposed are capable of being successively viewed while being projected onto a screen.\nTowards the end of achieving the desired and critical fluid thickness on the film strip, as well as for avoiding blemishes in the developed transparent images, the cassette contains a spring-biased pressure pad yieldably supporting the film strip between the pad and a fluid processor nozzle structure.\nThe current state-of-the art with respect to achieving this uniform and substantially blemish free coating is represented by the disclosure of U.S. Pat. No. 3,951,530 issued Apr. 20, 1976 to Frank M. Czumak, Paul B. Mason and Joseph A. Stella, which patent is commonly assigned with the present invention. In the disclosure of this patent, the film's emulsion surface is yieldably urged into sliding engagement with the processor nozzle so that a uniform and predetermined gap exists between the emulsion surface and a doctoring surface formed on the nozzle. Formation of the desired fluid thickness is effected after the fluid has been deposited onto the film and passes the doctoring surface positioned downstream, in the direction of film advancement during processing, of a nozzle opening. This formation is facilitated by the doctoring surface being configured to develop positive hydrodynamic pressures in the deposited fluid which force the film strip into sliding engagement with the pressure pad. To resist this tendency, the pressure pad spring must be constructed to retain the desired gap spacing despite the presence of these hydrodynamic pressures. Also, the biasing force of the spring is selected to provide a net balance of forces on the film strip. This is accomplished by selecting the spring biasing force to be substantially equal and opposite to the hydrodynamic pressures. As a result foreign particles, such as dust or the like, on the emulsion surface of the film will effect a slight instantaneous increase in the gap between the doctoring surface and such emulsion surface and be permitted to pass therebeneath. Consequently, foreign particles will not be trapped and the possibility of blemishes occurring on the developed film is substantially eliminated, it being understood that these particles, in themselves, are inconsequential to either processing or developing because of their minute size. It will be appreciated, therefore, the spring force of the spring must be carefully controlled, otherwise the foregoing functions of the spring will result in significant operational problems. The type of spring described in the foregoing identified patent, which is the kind typically used in cassettes of this category, is a thin leaf spring having one end staked to the cassette, and the other end contacting and biasing the pressure pad. While the leaf spring performs satisfactorily there is potential for problems arising from the fact that the spring can be overstressed and thereby have adversely altered the noted preselected biasing force. Usually the problems of overstressing arise by virtue of the cassette assembly operations wherein an operator must periodically bend and thereby possibly overstress the spring in order to produce the completed film cassette. Aside from this potential problem use of the conventional leaf spring in its present embodiment necessitates relatively costly assembly procedures because of the staking operation.\nIt follows, therefore, that avoidance of overstressing the spring biasing the pressure pad during assembly is a major focal point of attention to overall film cassette development.\nMoreover, the potential of the foregoing problems is further compounded by the requirement the cassette and its components including the pressure pad spring must be capable of mass production manufacturing techniques as well as be within tolerance levels incident to such techniques for the system to be acceptable in a competitive commercial market. Accordingly, the structural organization of components by which the film is supported in predetermined relationship to the doctoring surface is important to the overall system."}
{"text": "1. Field of the Invention\nThis invention is related to laser systems. In particular, this invention deals with adjusting alignment of laser beams in laser systems.\n2. Description of the Related Art\nWhen fabricating memory circuits, a laser repair system can be used to selectively sever conductive links, effectively removing faulty memory cells from the circuit.\nAs the size and spacing of link elements decreases, laser repair systems have had to increase in accuracy in order to perform their intended function. The complexity of a laser repair system capable of such accurate operation is significant. Multiple mirrors and other optical elements are used to generate and position a laser beam spot for severing a conductive link. Like the circuit fabrication process itself, laser repair systems are subject to many complex factors. For example, thermal expansion may lead to changes in the orientation or position of optical elements in the path of a laser beam. These changes to the elements that affect the laser beam can cause the laser beam spot to drift away from its intended location and can cause errors when trying to repair a circuit. Although the beam spot position is aligned with reference to wafer alignment markers with every new wafer processed, a misaligned laser beam path that deviates from a normal orientation to the work surface can still produce beam spots of unintended location, shape and/or size which adversely affect operation of the repair system.\nU.S. Pat. No. 6,483,071 (hereinafter referred to as the '071 patent) entitled “Method and system for precisely positioning a waist of a material-processing laser beam to process microstructures within a laser-processing site” is assigned the assignee of the present invention. The disclosure of the '071 patent is hereby expressly incorporated by reference in its entirety. The '071 patent discloses many features of a laser based system for memory repair, and is particularly related to accurate (sub-micron) and high-speed positioning of a laser beam waist relative to a link or similar target structure. In the '071 patent, an air-bearing based assembly was disclosed for positioning of optical components (e.g: an objective lens) along the optical (Z) axis. In addition to noise and reliability issues (ie: wearing mechanical parts) it was recognized that X,Y displacement errors during Z axis motion are much better controlled or eliminated with an air bearing system. Such displacements, even if a fraction of a micron, can lead to link severing results which are incomplete (e.g. contamination) or possibly cause damage to surrounding structures. Hence, a displacement of a laser beam from a target location by a fraction of one-micron, corresponding to a fraction of one spot diameter, may generally lead to reduced yield.\nTraditionally, laser repair systems have undergone periodic, manual adjustment to correct problems with alignment. For example, every month, a trained technician may have to manually adjust optical elements in order to correct alignment problems that have developed since the last adjustment. In the M430 laser link blowing machine from GSI, coarse adjustments to laser beam alignment were made by manually adjusting the laser beam orientation while viewing the laser beam spot with a “thru-lens viewing system” (TTLV). The TTLV is essentially a camera and TV monitor arrangement coupled to the laser beam path. The spot position was determined relative to a crosshair. The beam was first aligned to be centered in the lens aperture. Then the beam was aligned for zero spot translation during zoom expansion. Zoom adjustments corresponded to a range of spot sizes. If the beam was properly aligned along the Z-axis, the beam would appear stationary on the monitor for all zoom settings. Finer beam alignment was carried out by adjusting the spot size to a minimum, placing a calibration grid on the work surface, and performing iterative manual adjustments of turning mirrors to align the optical system and reduce any lateral (X-Y) displacement to within a specified tolerance.\nThis traditional approach to adjusting the alignment of a laser beam has several drawbacks. For example, the means used by the technician to determine beam alignment may itself be subject to error. Alignments based on erroneous alignment data may augment alignment problems in the system. Other problems may include the significant time expense involved in manual adjustment. Delays arising from manual alignment can represent a serious cost for businesses operating laser repair systems. For these reasons and others, automated methods of static laser beam alignment have been developed. Such methods are described for example in U.S. Pat. Nos. 5,011,282 to Ream, et al., 5,315,111 to Burns, et al., 5,923,418 to Clark et al., and 6,448,999 to Utterback et al. Of these prior patents, Burns, Clark, and Utterback split off portions of the laser beam to optical detectors placed adjacent to the laser beam path. Alignment of the beam with respect to the detectors is used to deduce alignment of the beam to the workpiece. In the '282 patent to Ream, changes in laser beam spot position on a target are used to determine a laser beam deviation angle, which can then be used to correct the laser beam path alignment."}
{"text": "Ocular hypotensive agents are useful in the treatment of a number of various ocular hypertensive conditions, such as post-surgical and post-laser trabeculectomy ocular hypertensive episodes, glaucoma, and as presurgical adjuncts.\nGlaucoma is a disease of the eye characterized by increased intraocular pressure. On the basis of its etiology, glaucoma has been classified as primary or secondary. For example, primary glaucoma in adults (congenital glaucoma) may be either open-angle or acute or chronic angle-closure. Secondary glaucoma results from pre-existing ocular diseases such as uveitis, intraocular tumor or an enlarged cataract.\nThe underlying causes of primary glaucoma are not yet known. The increased intraocular tension is due to the obstruction of aqueous humor outflow. In chronic open-angle glaucoma, the anterior chamber and its anatomic structures appear normal, but drainage of the aqueous humor is impeded. In acute or chronic angle-closure glaucoma, the anterior chamber is shallow, the filtration angle is narrowed, and the iris may obstruct the trabecular meshwork at the entrance of the canal of Schlemm. Dilation of the pupil may push the root of the iris forward against the angle, and may produce pupillary block and thus precipitate an acute attack. Eyes with narrow anterior chamber angles are predisposed to acute angle-closure glaucoma attacks of various degrees of severity.\nSecondary glaucoma is caused by any interference with the flow of aqueous humor from the posterior chamber into the anterior chamber and subsequently, into the canal of Schlemm. Inflammatory disease of the anterior segment may prevent aqueous escape by causing complete posterior synechia in iris bombe and may plug the drainage channel with exudates. Other common causes are intraocular tumors, enlarged cataracts, central retinal vein occlusion, trauma to the eye, operative procedures and intraocular hemorrhage.\nConsidering all types together, glaucoma occurs in about 2% of all persons over the age of 40 and may be asymptotic for years before progressing to rapid loss of vision. In cases where surgery is not indicated, topical .beta.-adrenoreceptor antagonists have traditionally been the drugs of choice for treating glaucoma.\nIt has long been thereof that one of the sequelae of glaucoma is damage to the optic nerve head. This damage, referred to as \"cupping\", results in depressions in areas of the nerve fiber of the optic disk. Loss of sight from this cupping is progressive and can lead to blindness if the condition is not treated effectively.\nUnfortunately lowering. intraocular pressure by administration of drugs or by surgery to facilitate outflow of the aqueous humor is not always effective in obviating damage to the nerves in glaucomatous conditions. This apparent contradiction is addressed by Cioffi and Van Buskirk [Surv. of Ophthalmol., 38, Suppl. p. S107-16, discussion S116-17, May 1994] in the article, \"Microvasculature of the Anterior Optic Nerve\". The abstract states:\nThe traditional definition of glaucoma as a disorder of increased intraocular pressure (IOP) oversimplifies the clinical situation. Some glaucoma patients never have higher than normal IOP and others continue to develop optic nerve damage despite maximal lowering of IOP. Another possible factor in the etiology of glaucoma may be regulation of the regional microvasculature of the anterior optic nerve. One reason to believe that microvascular factors are important is that many microvascular diseases are associated with glaucomatous optic neuropathy.\nSubsequent to Cioffi, et al., Matusi published a paper on the \"Ophthalmologic aspects of Systemic Vasculitis\" [Nippon Rinsho, 52 (8), p. 2158-63, August 1994] and added further support to the assertion that many microvascular diseases are associated with glaucomatous optic neuropathy. The summary states:\nOcular findings of systemic vasculitis, such as polyarteritis nodosa, giant cell angitis and aortitis syndrome were reviewed. Systemic lupus erythematosus is not categorized as systemic vasculitis, however its ocular findings are microangiopathic. Therefore, review of its ocular findings was included in this paper. The most common fundus finding in these diseases is ischemic optic neuropathy or retinal vascular occlusions. Therefore several points in diagnosis or pathogenesis of optic neuropathy and retinal and choroidal vaso-occlusion were discussed. Choroidal ischemia has come to be able to be diagnosed clinically, since fluorescein angiography was applied in these lesions. When choroidal arteries are occluded, overlying retinal pigment epithelium is damaged. This causes disruption of barrier function of the epithelium and allows fluid from choroidal vasculatures to pass into subsensory retinal spaces. This is a pathogenesis of serous detachment of the retina. The retinal arterial occlusion formed non-perfused retina. Such hypoxic retina released angiogenesis factors which stimulate retinal and iris neovascularizations and iris neovascularizations may cause neovascular glaucoma.\nB. Schwartz, in \"Circulatory Defects of the Optic Disk and Retina in Ocular Hypertension and High Pressure Open-Angle Glaucoma\" [Surv. Ophthalmol., 38, Suppl. pp. S23-24, May 1994] discusses the measurement of progressive defects in the optic nerve and retina associated with the progression of glaucoma. He states:\nFluorescein defects are significantly correlated with visual field loss and retinal nerve fiber layer loss. The second circulatory defect is a decrease of flow of fluorescein in the retinal vessels, especially the retinal veins, so that the greater the age, diastolic blood pressure, ocular pressure and visual field loss the less the flow. Both the optic disk and retinal circulation defects occur in untreated ocular hypertensive eyes. These observations indicate that circulatory defects in the optic disk and retina occur in ocular hypertension and open-angle glaucoma and increase with the progression of the disease.\nThus, it is evident that there is an unmet need for agents that have neuroprotective effects in the eye that can stop or retard the progressive damage that occurs to the nerves as a result of glaucoma or other ocular afflictions.\nProstaglandins were earlier regarded as potent ocular hypertensives; however, evidence accumulated in the last two decades shows that some prostaglandins are highly effective ocular hypotensive agents and are ideally suited for the long-term medical management of glaucoma. (See, for example, Starr, M. S. Exp. Eye Res. 1971, 11, pp. 170-177; Bito, L. Z Biological Protection with Prostaglandins Cohen, M. M., ed., Boca Raton, Fla., CRC Press Inc., 1985, pp. 231-252; and Bito, L. Z., Applied Pharmacology in the Medical Treatment of Glaucomas Drance, S. M. and Neufeld, A. H. eds., New York, Grune & Stratton, 1984, pp. 477-505). Such prostaglandins include PGF.sub.2.alpha., PGF.sub.1.alpha., PGE.sub.2, and certain lipid-soluble esters, such as C.sub.1 to C.sub.5 alkyl esters, e.g. 1-isopropyl ester, of such compounds.\nIn the U.S. Pat. No. 4,599,353 certain prostaglandins, in particular PGE.sub.2 and PGF.sub.2.alpha. and the C.sub.1 to C.sub.5 alkyl esters of the latter compound, were reported to possess ocular hypotensive activity and were recommended for use in glaucoma management.\nAlthough the precise mechanism is not yet known, recent experimental results indicate that the prostaglandin-induced reduction in intraocular pressure results from increased uveoscleral outflow [Nilsson et al., Invest. Ophthalmol. Vis. Sci. 28 (suppl), 284 (1987)].\nThe isopropyl ester of PGF.sub.2.alpha. has been shown to have significantly greater hypotensive potency than the parent compound, which was attributed to its more effective penetration through the cornea. In 1987 this compound was described as \"the most potent ocular hypotensive agent ever reported.\" [See, for example, Bito, L. Z., Arch. Ophthalmol. 105, 1036 (1987), and Siebold et al., Prodrug 5, 3 (1989)].\nWhereas prostaglandins appear to be devoid of significant intraocular side effects, ocular surface (conjunctival) hyperemia and foreign-body sensation have been consistently associated with the topical ocular use of such compounds, in particular PGF.sub.2.alpha. and its prodrugs, e.g. its 1-isopropyl ester, in humans. The clinical potential of prostaglandins in the management of conditions associated with increased ocular pressure, e.g. glaucoma, is greatly limited by these side effects.\nCertain phenyl and phenoxy mono, tri and tetra nor prostaglandins and their 1-esters are disclosed in European Patent Application 0,364,417 as useful in the treatment of glaucoma or ocular hypertension.\nIn a series of co-pending United States patent applications assigned to Allergan, Inc. prostaglandin esters with increased ocular hypotensive activity accompanied by no or substantially reduced side-effects are disclosed. The co-pending U.S. Ser. No. 386,835 (filed Jul. 27, 1989), relates to certain 11-acyl-prostaglandins, such as 11-pivaloyl, 11-acetyl, 11-isobutyryl, 11-valeryl, and 11-isovaleryl PGF.sub.2.alpha.. Intraocular pressure reducing 15-acyl prostaglandins are disclosed in the co-pending application U.S. Ser. No. 357,394 (filed May 25, 1989). Similarly, 11,15- 9,15- and 9,11-diesters of prostaglandins, for example 11,15-dipivaloyl PGF.sub.2.alpha. are known to have ocular hypotensive activity. See the co-pending patent applications U.S. Ser. No. 385,645 filed Jul. 27, 1990, now U.S. Pat. Nos. 4,494,274; 584,370 which is a continuation of U.S. Ser. Nos. 386,312, and 585,284, now U.S. Pat. No. 5,034,413 which is a continuation of U.S. Ser. No. 386,834, where the parent applications were filed on Jul. 27, 1989. The disclosures of these patent applications are hereby expressly incorporated by reference.\nFinally, certain EP.sub.2 -receptor agonists are disclosed in Nials et al, Cardiovascular Drug Reviews, Vol. 11, No. 2, pp. 165-179, Coleman et al, Comprehensive Medicinal Chemistry, Vol. 3, pp. 643-714, 1990 and Woodward et al, Prostaglandins, pp. 371-383, 1993."}
{"text": "A system for integrating a third generation (3G) cellular network with a WLAN has been proposed by the 3rd Generation Partnership Project (3GPP) in the standards document 3GPP TS 23.234, entitled “3GPP system to Wireless Local Area Network (WLAN) Interworking; System Description (Release 6),” which is incorporated herein by reference."}
{"text": "Example computer applications are the so-called “business applications.” A business application may be any software or set of computer programs that are used by business users to perform various business functions by processing and analyzing data stored in databases. Available business applications relate to different business functions including, for example, customer relationship management, enterprise asset management, enterprise resource planning, financials, human capital management, procurement, product lifecycle management, supply chain management, and sustainability, etc. Other business applications may relate to functions such as business intelligence, data warehousing, enterprise information management, enterprise performance management, governance, risk, and compliance.\nBusiness applications may be interactive. A business application may provide a user interface (UI) through which a user can query data and view query results. A user's queries may be directed to diverse or large databases (e.g., business information warehouses). Use of some business applications (e.g., business analytics applications) may involve numerous queries and extensive query data analysis and processing. As today's computing environments evolve, business applications are becoming more and more complex. In a distributed environment, business applications critical to an enterprise may be spread across multiple systems and access multiple databases. At least in large system implementations, query result display processes may be relatively slow.\nConsideration is now being given to aspects of displaying query results to users."}
{"text": "An exhaust bypass valve is often used to control operation of serial turbocharger systems. Such a valve may be operated to physically divert exhaust or alter pressures in exhaust pathways, for example, to direct exhaust flow partially or fully to one of multiple turbines in a system. During operation, a typical exhaust bypass valve experiences high exhaust pressure on one side and lower pressure on the other side. To effectively seal the high pressure environment from the low pressure environment, considerable force is required to maintain contact between a valve and a valve seat. In a sealed state of a valve and valve seat, pressure differentials may challenge one or more inter-component seals and result in detrimental exhaust leakage. Various technologies described herein have potential to reduce cost as well as provide for effective exhaust bypass valve sealing."}
{"text": "The block diagram of FIG. 1 shows one example of a prior art cordless telephone. A battery 1, which is chargeable through a charging terminal 2, supplies power to a power source circuit 3 by closing a power switch 4. Power source circuit 3, under control of a control circuit 5, converts the voltage from battery 1 to an appropriate voltage, and supplies this voltage to control circuit 5, to a transmitter 6 and to a receiver 7.\nIn a waiting state, i.e., awaiting an incoming call, the converted voltage from power source circuit 3, controlled by circuit 5, is intermittently applied to receiver 7. Upon receipt of a group signal, the receiver is synchronized with the transmitted signal from a base station for receipt of further intermittent signals. Consequently, receiver 7 carries out reception operation intermittently as shown in FIG. 2, thereby saving power consumption of the cordless telephone.\nA significant problem in these prior art systems is the inability of the user of the battery powered radio to detect whether the system is operating properly. In particular, it has become necessary for the user to determine whether the communication link from the base station to the radio receiver is operating properly. Any defect in the base transmitter, or in the radio receiver, or due to interference could prevent the reception and processing of the data at the receiver. Not only is it important to identify the existence of a problem but to indicate it to the user in a manner which will avoid a large consumption of power.\nThe synchronized operation of several battery powered radio devices incorporating the prior art battery saving function is shown in FIGS. 3(a)-(c).\nReferring to FIG. 3(b), all devices in a system are divided into M groups, for example, 3 groups. A base unit repeatedly broadcasts a frame signal consisting of M group signals. Each group signal includes preamble words and N calling words (FIG. 3(a)). The preamble word includes a predetermined synchronization signal for synchronization of received signals and a group identification signal for showing that the following paging words are addressed to pagers belonging to the group designated by the group identification signal. In response to a calling request, an identification signal of the device to be called is assigned to one of calling words 1 to N.\nReferring to FIG. 3(c), if a device is turned on, the device examines the received signals to detect the synchronization signal. Once detecting the synchronization signal, the device performs intermittent reception so that only a group signal to the device is received.\nIn fact, prior art radio devices such as the cordless devices mentioned above have not even provided any indication whether the power source switch was closed. Apparently, if such an indicator was used to display the activation of the source, the display would emit light continuously thereby increasing power consumption and shortening the active life of the radio devices. It would be possible to increase the capacity of the battery, but this would necessarily result in high cost, increased weight and longer charging times."}
{"text": "In developments of emerging technology, such as new wireless standards or Artificial Intelligence, the amount of data required to be processed is increasing substantially. With the profusion of data, more computational requirements are placed on general purpose CPUs, specialized CPUs (i.e. GPU, TPU) and/or specialized Hardware Accelerators to expeditiously process the data.\nAs the computational requirement placed on the processors increases, the performance of the processors is often inadequate to handle computationally intensive tasks on large amount of data. In some cases, even if specialized processors are capable of handling the computational requirements, the cost of such processors is often prohibitive for many applications.\nThere are various factors which limit the computational capabilities of a processor. Traditionally, the processors use internal registers to temporarily hold the source input data which are loaded from the data memory. The processor then performs an arithmetic or other programmed operation using the values stored in the temporary registers as the operands, and writes the result of the operation to another temporary register. Finally, the processor stores the result in the temporary register back to the data memory.\nFor performing such operations, many instructions are required. For example, ADD Immediate instructions to calculate the operand addresses; LOAD Instructions to load the operands; MULTIPLY instruction to multiply the operands; ADD Immediate instruction to calculate the destination address; and STORE instruction to write the result to the destination memory location.\nDuring the execution of these instructions, due to the inherent load/store latency associated with the data memory and the limited availability of the temporary registers, the instruction executions are often blocked by pipeline stalls resulting in degraded processor performance. The problem of pipeline stall is compounded when the processor operates on large sets of data.\nOther common techniques employed in the industry, such as SIMD and Vector Instruction Extensions, try to address the performance issue by parallel data processing. However, these techniques, even though they obtain performance increase through parallelism, are still subject to the aforementioned limitations.\nTherefore, it is desirable to have a flexible solution capable of processing large amount of data, which can also be quickly programmed, deployed, and modified as the product matures."}
{"text": "This invention concerns mechanical connections between control rod ends and a pivot pin, as are commonly used to connect control rods to a pin in transmission linkages for automobiles.\nIn these connections, it is now usual to incorporate a soft elastomeric isolator between the rod and the pin to minimize the transmission of vibrations to a shift lever for example which is gripped by the hand of the driver of an automobile. A disadvantage of such isolators has been the looseness or lash introduced into the connection by the pin contact compressing the isolator material.\nSuch looseness has also been introduced by the clearances necessary to assemble the component parts of the connection.\nLabor costs are also incurred by the need to fit together complex components to complete the connection when the automobile is assembled, and the need to insure proper assembly.\nAnother difficulty has been involved in attempting to reduce the force necessary to fit the pin to the connection to simplify assembly while still insuring that an adequately high separation force would be required to disconnect the pin after the connection is made to prevent unintended disconnection.\nIt is the object of the present invention to provide an arrangement and method for connecting one end of a rod to a headed pin which incorporates a vibration isolator without introducing excessive looseness in the connection.\nIt is another object of the present invention to provide such an arrangement and method in which there are minimal clearances necessary for assembly purposes.\nIt is yet another object of the present invention to provide such an arrangement and method which requires only low installation forces while insuring that separation forces are sufficiently high to preclude unintended disconnection of the components."}
{"text": "U.S. Pat. No. 2,070,263 describes a method for obtaining aqueous solutions of hydrobromic acid which consists, in a first step, in passing hydrogen through liquid bromine maintained at a temperature of between 37.degree. C. and 42.degree. C. in order to form a mixture of bromine and hydrogen gas which is burnt at a temperature of between 600.degree. C. and 850.degree. C. By working in this way, it is difficult to have an intimate mixture of bromine and hydrogen in stoichiometric amounts on account of the difficulties in rigorously maintaining the temperatures and the thermodynamic equilibria.\nAn instability of the combustion flame has also been observed, in devices based on the direct combustion of bromine in hydrogen according to an H.sub.2 /Br.sub.2 molar ratio of greater than 1, and for which no technique for mixing the reactants is mentioned, this instability being manifested in particular by strong vacillation of the flame at the burner outlet, going as far as a detachment of the flame (\"blow off\") from the said burner, which may entail a risk of explosion and a fluctuating quality of the hydrogen bromide gas produced.\nFurthermore, such flames become extended forming cones at the base of which are regions from which the reactants are liable to escape without being burnt.\nThis disrupts the combustion of the bromine in the hydrogen and results especially in residual bromine in the combustion gases, this being of a nature to bring about a considerable decrease in the lifetime of the burners, limit the range of materials which may be used and degrade the quality of the hydrogen bromide gas, thereby preventing it from being used as a reactant for downstream syntheses (secondary reactions, colorations of the products) or for the preparation of hydrobromic acid solutions.\nPatent FR 2,365,516 proposes a process which improves the stability of the flame resulting from the combustion of bromine in hydrogen by establishing a helical stream of bromine in a cylindrical chamber, then injecting the hydrogen radially towards the outside in the helical stream of bromine and continuously supplying a flame close to the chamber with the helical stream of bromine and hydrogen.\nThis process, using a molar excess of hydrogen of 2.6%, leads to an HBr gas containing 300 ppm of bromine by volume, which still gives rise to colorations of the downstream synthesis products as well as the drawbacks mentioned above.\nIn addition, the complexity of the burner entails a lack of flexibility. Thus, in particular, when it is desired to increase the capacity of the said device, several burners are arranged side-by-side in the same chamber. In such an arrangement, it cannot be avoided that the flames from different burners mounted in parallel will interfere with each other, and furthermore this arrangement is unacceptable with regard to obtaining good distribution of the reactants. This configuration inevitably leads to a lowering in the degree of conversion of the bromine, complicates the control of the cooling of the HBr formed and increases the risks of explosion."}
{"text": "This invention relates in general to fire resistant structures and, more particularly, to a fire resistant door.\nFire barriers are included within the design of many types of buildings in order to block the spread of a fire once it has been ignited within the building. Because the placement of door openings in the walls of the building provides an avenue for the fire to spread from room to room, much attention has been focused on the designing of doors which are fire resistant and can impede the spread of the fire.\nPanel doors have a plurality of flat or raised panels interconnected with vertically extending stiles and horizontal rails. Although panel doors are widely utilized because of their visually pleasing appearance, they generally have poor resistance to fire because air is able to infiltrate the door at the juncture of the panels with the stiles and rails. The fire is fed by the oxygen present in the air which seeps through the panel joints and can quickly burn through the door at those joints. In an effort to increase their fire resistance, some panel doors are available which have segments of fire resistant material inserted within grooves milled into the edges of the stiles, rails and panels at the junctures of those components. The placement of the fire resistant material at those locations is generally effective to slow the rate at which the fire can burn through the panel joints, but the fire is still able to burn through the wooden portions of the door at a faster rate than is desired in many instances.\nAnother type of conventional fire resistant door mimics the appearance of a panel door by applying half panels to a core of fire resistant material. The core comprises a solid sheet of material which is milled on both faces to form recesses at the intended location of the panels. Half panels are then positioned within the recesses and veneer and trim are applied to the exposed surfaces of the core. The resulting door can be very resistant to fire because the solid core blocks any air infiltration through the door. The milling operation, however, may reduce the structural integrity of the door and result in warping, sagging or other deformation of the door.\nA need thus exists for a panel door which has an enhanced fire resistance but maintains the desired structural and visual appearance."}
{"text": "This invention is generally directed to layered photoresponsive imaging devices, and more specifically to photoconductive devices having incorporated therein certain novel bisazo compounds. Therefore, in one embodiment of the present invention there are provided photoconductive layered imaging members comprised of certain bisazo compounds and arylamine hole transport layers. In one important embodiment of the present invention, there is provided a photoresponsive device comprised of various specific bisazo compounds, including 4,4'-bis(1\"-azo-2\"-hydroxy-3\"-naphthanilide)-1,1'-dianthraquinonylamine; 4,4'-bis(1\"-azo-2\"-hydroxy-3\"-naphtho-p-trifluoromethylanilide-1,1'-dianth raquinonylamine; and the derivatives thereof; and wherein the member further includes therein a charge, or hole transport layer. The aforementioned photoconductors possess a number of advantages indicated hereinafter inclusive of high photosensitivity, excellent photosensitivity to wavelengths of from about 400 to 750 nanometers, and high cyclic stability; and further, are very economical enabling these devices to be readily disposable. Accordingly, the photoresponsive imaging members of the present invention are useful in various electrophotographic and electrostatographic imaging processes wherein, for example, latent images are formed thereon followed by development and transfer to a suitable substrate. More specifically, the imaging members of the present invention with photosensitivity of from about 400 to about 750 nanometers enabled such members to be useful for electrophotographic imaging devices, and also these members can be incorporated into light emitting diode printers as well as multifunctional printer electrophotographic apparatuses.\nNumerous different xerographic photoconductive members are known including, for example, a homogeneous layer of a single material such as vitreous selenium, or a composite layered device containing a dispersion of a photoconductive composition. An example of one type of composite xerographic photoconductive member is described, for example, in U.S. Pat. No. 3,121,006 wherein there is disclosed finely dispersed divided particles of a photoconductive inorganic compound dispersed in an electrically insulating organic resin binder. These members contain, for example, coated on a paper backing, a binder layer containing particles of zinc oxide uniformly dispersed therein. The binder materials disclosed in this patent comprise a material such as polycarbonate resins, polyester resins, polyamide resins, and the like, which are incapable of transporting for any significant distance injected charge carriers generated by the photoconductive particles.\nThere are also known photoconductive members comprised of inorganic or organic materials wherein the charge carrier generating, and charge carrier transport functions are accomplished by discrete continuous layers. Additionally, layered photoconductive members are disclosed in the prior art which include an overcoating layer of an electrically insulating polymeric material.\nRecently, there have been disclosed other layered photoresponsive devices including those comprised of separate generating layers, and transport layers as described in U.S. Pat. No. 4,265,990, and overcoated photoresponsive materials containing a hole injecting layer overcoated with a hole transport layer, followed by an overcoating of a photogenerating layer, and a top coating of an insulating organic resin, reference U.S. Pat. No. 4,251,612. Examples of photogenerating layers disclosed in these patents include trigonal selenium and phthalocyanines, while examples of transport layers include certain diamines as mentioned herein. The disclosures of each of these patents, namely U.S. Pat. Nos. 4,265,990 and 4,251,612, are totally incorporated herein by reference.\nMany other patents are in existence describing photoresponsive devices including layered devices containing generating substances, such as U.S. Pat. No. 3,041,167 which discloses an overcoated imaging member containing a conductive substrate, a photoconductive layer, and an overcoating layer of an electrically insulating polymeric material. This member is utilized in an electrophotographic copying system by, for example, initially charging the member with an electrostatic charge of a first polarity, and imagewise exposing to form an electrostatic latent image which can be subsequently developed to form a visible image.\nFurthermore, there are disclosed in U.S. Pat. Nos. 4,232,102 and 4,233,383 photoresponsive imaging members comprised of trigonal selenium doped with sodium carbonate, sodium selenite, and trigonal selenium doped with barium carbonate, and barium selenite or mixtures thereof. Moreover, there are disclosed in U.S. Pat. No. 3,824,099 certain photosensitive hydroxy squaraine compositions. According to the disclosure of this patent, the squaraine compositions are photosensitive in normal electrostatographic imaging systems.\nAlso known are photoconductive members containing therein various squaraine compositions. Thus, for example, there are illustrated in U.S. Pat. No. 4,508,803, the disclosure of which is totally incorporated herein by reference, photoconductive devices containing novel benzyl fluorinated squaraine compositions. Specifically, in one embodiment illustrated in the '803 patent there is described an improved photoresponsive device comprised of a supporting substrate, a hole blocking layer, an optional adhesive interface layer, an inorganic photogenerating layer, a photoconducting composition layer comprised of benzyl fluorinated squaraine compositions, and a hole transport layer. Other representative patents disclosing photoconductive devices with squaraine components therein, or processes for the preparation of squaraines include U.S. Pat. No. 4,507,408; 4,552,822; 4,559,286; 4,507,480; 4,524,220; 4,524,219; 4,524,218; 4,525,592; 4,559,286; 4,415,639; 4,471,041; and 4,486,520. The disclosures of each of the aforementioned patents are totally incorporated herein by reference.\nFurther, disclosed in the prior art are composite electrophotographic photosensitive materials with various bisazo compounds. For example, there are illustrated in Japanese Ricoh Patent Publication No. 6064354, published Apr. 12, 1985, composite photoconductors wherein one of the photoconductor layers contain a bisazo compound of the formulas as illustrated. Further, there are illustrated in several U.S. patents layered organic electrophotographic photoconductor elements with bisazo, trisazo, or related compounds. Examples of these U.S. pat. Nos. include 4,596,754; 4,555,567; 4,555,667; 4,440,845; 4,486,522; 4,486,800; 4,299,896; 4,551,404; 4,309,611; 4,418,133; 4,293,628; 4,427,753; 4,495,264; 4,359,513; 3,898,084; 4,400,455; 4,390,608; 4,327,168; 4,299,896; 4,314,015; 4,486,522; 4,486,519; and Konishiroku Japanese Patent Laid Open Publication No. 60111247.\nAlso of interest is U.S. Pat. No. 4,713,307, which illustrates photoconductive imaging members containing a supporting substrate, certain azo pigments of 2,7-bis(1'-azo-2'-hydroxy-3'-naphthanilide) naphthalene, and the derivatives thereof; and a charge transport layer.\nAlthough photoconductive imaging members with bisazo compounds are known, there remains a need for novel bisazo photoconductor devices with extended red or near-IR photoresponses thereby enabling their selection in imaging apparatus with light emitting diodes. Additionally, there continues to be a need for layered photoresponsive imaging members have incorporated therein certain bisazo compounds, which members will enable the generation of acceptable high quality images, and wherein these members can be repeatedly used in a number of imaging cycles without deterioration thereof from the machine environment or surrounding conditions. Moreover, there is a need for improved layered photoresponsive imaging members wherein the bisazo compounds selected for one of the layers are substantially inert to the users of such members. Additionally, there is an important need for layered photoconductors with bisazo compounds, which photoconductors are of low cost, high sensitivity, and possess high cyclic stability. There also is a need for bisazo photoconductors that possess photosensitivity in the wavelength region of from about 650 to about 750 nanometers enabling these photoconductors to be selected for electrophotographic, particularly xerographic, imaging processes; light emitting diode printers; and multifunctional printer electrophotographic apparatuses."}
{"text": "Telephony has advanced dramatically with the advancement of technology. Telephone communication was once limited to an analog public switched telephone network (PSTN), where the PSTN has been traditionally formed of two types of telephone carriers, local and long distance. The local carriers established local networks for subscribers to communicate within local regions, and the long distance carriers created networks between the local networks to enable subscribers of different local carriers to communicate with one another.\nOver time, mobile telephone networks were developed to enable subscribers to use mobile telephones. At first, the wireless networks and handsets were analog. Technology for the wireless networks was developed to provide digital wireless communications, which provided a clearer signal than analog wireless communications.\nAbout the same time that the digital wireless networks were developing, the Internet was also developing. The International Standards Organization (ISO) developed an Open Systems Interconnection (OSI) basic reference model in 1977 that currently includes seven different layers. Each of the layers provides protocols for certain types of operations. More specifically, the seven layers include: physical layer (Layer 1), data link layer (Layer 2), network layer (Layer 3), transport layer (Layer 4), session layer (Layer 5), presentation layer (Layer 6), and application layer (Layer 7). Each entity interacts directly with the layer immediately beneath it and provides facilities for use by the layer above it. The protocols on each layer enable entities to communicate with other entities on the same layer. The Internet initially provided for simple digital data to be communicated between users. One of the early communication application included email. However, as communications standards and protocols developed, the Internet matured to include more advanced communication applications, including voice over Internet protocol (VoIP).\nFIG. 1 is an illustration of a legacy telecommunications network that includes class 4 and 5 switches 102a-102n (collectively 102) and 104a-104n (collectively 104), respectively, connected to a signaling system #7 (SS7) network 106 (indicated as dashed lines). Historically, the class 5 switches 104 were generally configured to communicate via in-band signaling verses the use of SS7 signaling and operate to form a local network of subscribers within the network to place telephone calls to one another. The class 4 switches 102 were developed for long distance connections between the class 5 switches 104 at end offices (not shown). The class 4 switches 102, which are monolithic, are generally formed of multiple components, including a port, port cross-connect matrix, switch messaging bus with external signaling units, and call processing unit, as understood in the art. Class 4 switches are circuit based and utilize time division multiplexing (TDM) and are capable of terminating higher high-speed communications, including T1, T3, OC-3, and other four-wire circuit connections. As understood in the art, TDM is a synchronous communications protocol.\nThe SS7 network 106, which includes signal transfer points (STPs) 108a-108n (collectively 108), service switching points (SSs) on the class 4 and class 5 switches, and service control points (SCPs—not shown). The SS7 network is connected to the class 4 and 5 switches for providing and maintaining inter-switch call services between the switches. The SS7 network is used to signal out-of-band call setup, as out-of-band signaling is more secure and faster than in-band signaling. The call state changes of the inter-switch trunks of the class 4 switch are communicated to the adjoining switches via the SS7 network via a connection to the STP. To manage and route calls, the STPs 108 are used as an Inter-Switch messaging network whereby two switches control the trunking between the switches via messaging over the SS7 network provided by the STP switches that act as the inter-switch message bus. A call state machine of the class 4 switch provides control for routing traffic within the cross-connect matrix of the monolithic switch. The call state machine also provides call control signaling information to other switches via the connections to the STPs. The call control signaling information is routed via the STPs to other switches for call setup and tear-down. The call control information routed by the STPs contains pertinent information about the call to allow the terminating switch to complete various calls.\nTelephony has benefited from the development of the OSI model in a vast number of ways. One way has been through separating the call controller into a distributed cross-connect on an asynchronous network, such as asynchronous transfer mode (ATM or a Internet Protocol (IP) network. FIG. 2 is an illustration of a conventional telephony network 200 that includes a packet network 202. In one embodiment, the packet network 202 is an ATM network. Media gateways (MGs) 204a-204n (collectively 204) are media translation or conversion devices that modify and convert protocols between disparate communication networks. The media gateways 204, which are in communication with class 5 switches 206a-206n (collectively 206) are located at the edge of the packet network 202. The media gateways 204 convert TDM packets or streams 208 into packets, frames, or cells (collectively referred to hereinafter as “packets”) 210 and vice versa.\nThe packet network 202 operates independently as a distributed virtual media gateway port cross-connect for voice calls primarily due to one or more call control managers (CCMs) 212 located on the packet network 202. The call control manager 212 is in communication with the media gateways 204 and operates to control the media gateways 204 and provide instructions on how to rotate the packets 210 via far-end address allocations. By separating the call controller from the class 4 switches, the packet network 202 becomes, in effect, the virtual cross-connect of the switching system. The packet network 202 enables packets 210 that include voice data, commonly known as bearer packets, to be tagged with a destination address 214a and origination address 214b for enabling content data 214c to be properly routed from the origination media gateway 204b to the destination media gateway 204a. The media gateways 204 use of the packet network 202 is controlled by the CCM 212 and may communicate the packets 210 over the packet network 202 via IP addresses and virtual circuit (VC) or virtual path (VP) between the media gateways 204 to appropriately route the packets to the correct destination network node through the packet network 202. CCM 212 receives call state processing information from the media gateways 204 and signaling points, and processes the call state changes by using look up tables (not shown). The CCM 212 thereafter communicates packet addressing and state changes to the media gateways 204 to process the call.\nEthernet protocol was developed to provide for a computer network that enables multiple computers to share a common external inter-communication bus. Ethernet is generally used to provide for local area networks (LANs). Ethernet operates by communicating frames of data. While Ethernet operates well within a local environment (e.g., within a building) because Ethernet assumes that there is an known capacity of bandwidth associated with the bus standards set forth in the IEEE 802.3 standard that defines Ethernet. Ethernet is a shared environment, where co-utilization creates transmission errors called collisions. These collisions are detected by Ethernet cards in computers and a random re-transmission timer is used to avoid the next collision. Ethernet poses special problems for use in communications systems given it lacks dedicated bandwidth and time slots. The shared nature of an Ethernet network creates additional complexities in that the amount of available bandwidth can vary when used with wireless technologies.\nCommunication protocols transmitted over packet networks, such as ATM or IP networks, may utilize TDM based transmission facilities, which are synchronous as compared to Ethernet transmission facilities, which are asynchronous. Synchronous transmission protocols utilize a common clock and channel schema so that each device on the network operates synchronously with a dedicated path. Two types of “connection” state knowledge are present in a dedicated system, such as a TDM. Each channel has a dedicated amount of bandwidth and an error rate that is calculated from a common clock to determine path errors. The two types of connection state awareness functionality are provided by the channel itself and the common clock and data within a TDM header. The common clock provides for a determination of (i) a communications data rate from one end-point to another end-point and (ii) the data quality. Additionally, the TDM protocol includes “far end state” data in a header of a TDM frame to indicate whether there is a connection at the far end, thereby providing an indication of continuity along the communications path. Specifically, in-band end-to-end alarming allows the cross-connect devices to receive indications of continuity problems with other end-points. The in-band alarming is also provided for connection quality, where Bit Error Rate (BER) allows each end-point to know the quality of the data being received. Furthermore, bandwidth is always in use meaning that packets are synchronous, which that the far end knows exactly how many packets are to be sent and received in a given time period (e.g., one second). Computation of utilization is easily made by using the known bandwidth and multiplying it by the “seizure” time or the amount of time in use.\nPacket-based communications sessions lack circuit based connection state awareness indicators and clocking functionality to provide a session controller the ability to know the path connectivity state to efficiently manage making call handling decisions with anything other than ample bandwidth to setup and use sessions. This lack of connection path state awareness with the communications protocols, such as Ethernet and Internet protocol (IP) technologies, result in “gaps” in terms of being able to react to decaying transmission path quality and be sensitive to shared use of bandwidth. Most IP call controller solutions are founded on enterprise applications, where a single entity owns the network and scale of the network is relatively small. IP and Ethernet protocols lack the in-band path signaling, quality and use metrics to allow for this scale, or the ability to perform enhanced call handling with paths outside the governance of the packet network. Because packet communications are asynchronous, there is no common clock, and, thus, there is no way to know how many packets were transmitted, which, in turn, removes the ability to characterize transmission quality of the entire path, the amount of bandwidth available, or the amount in use. Further, packet networks are “converged” meaning they have both real-time and non-real-time bandwidth use. Currently, there is no in-band mechanism for determining real-time and non-real-time bandwidth use; having such information would allow for handling calls. It is commonly understood that proper connection operational assumptions are made by call control engines when the SS7 signaling path is properly operating (e.g., provisioned bandwidth is available) between end-points within the SS7 network. These operational assumptions are problematic as Ethernet, IP, and other data networks become oversubscribed and cause the packet network to become congested and prevent throughput. In cases where an end-point, such as a WiFi telephone, is mobile and bandwidth changes with signal strength (e.g., a WiFi telephone losing bandwidth as an individual walks away from a connection point antenna), the connection operation assumptions also fail to provide graceful call handling.\nOne available technique in packet networks to prevent oversubscription of real-time media traffic is through the use of call admission control (CAC) or the IP equivalent known as Resource Reservation Protocol (RSVP). CAC is primarily used to prevent congestion in voice traffic and is applied in the call setup phase to ensure there is enough bandwidth for data flow by reserving resources. To reserve bandwidth through the entire packet network, a CAC requires that the CAC procedure be performed at each point along a virtual circuit between two media gateways on which a call is to be routed, and often in a bi-directional fashion. While CAC functionality exists, the use of such CAC functionality is almost never applied because of the amount of time needed by the CAC procedure during call set up. For example, currently, CAC typically cannot operate over 40 calls per second and typical call set-ups on media gateways or class 4 switches may be 200 calls per second or higher.\nOne technique used to monitor the performance of IP session performance (i.e., after a call session has been established) is the use of the real-time control protocol (RTCP) as defined in IETF RFC 3550. RTCP collects statistics on a media connection, including bytes sent, packets sent, lost packets, jitter, feedback, and round trip delay. Other information may be provided in the RTCP packet using profile specific extensions. RTCP, which operates on a per session basis, is used for quality of service (QoS) reporting after termination of a session. The statistics information may be used, for example, to improve the quality of service by limiting data flow or changing CODEC compression. Utilization of the real-time QoS statistics, however, is limited to the specific session associated with the RTCP stream.\nAn emerging standard that is being developed for Ethernet performance measures is 802.1AG. This standard operates by generating and communicating an 802.1AG packet or “heart beat” over an Ethernet network segment. The 802.1AG packets are communicated via a Layer 2 Ethernet Virtual Circuit, such as a VLAN or Ethernet tunnel. At the ends and mid-points in Ethernet tunnels, 802.1AG packets are transmitted periodically over the Ethernet network to the far end. The Y.1731 protocol is utilized to calculate the number of data frames communicated between the 802.1AG packets. This configuration enables a performance measures (PM) to compute certain information about the performance of the path between the end-points on an Ethernet network. This combination of 802.1AG and Y.1731 enables the end points to be knowledgeable about the Frame Loss Rate (FLR), packet delay, and jitter in the path. This configuration is helpful to assist in monitoring performance of an Ethernet network path and diagnosing connectivity faults. However, the configuration falls short of providing the amount of real-time bandwidth in use or the total bandwidth in use. This information is useful to the proper management by a session controller handling calls during periods of flux in the packet transmission path, or the management of the real-time traffic.\nService providers often have trouble isolating and diagnosing network problems. To attempt to locate a packet loss problem along a node segment (i.e., a path between two network communications devices) over a network a probe that may be used to trace data packets being communicated over the node segment. This probe, however, is typically an external device from the network communications devices and operates to run a trace over an instant of time to determine network performance information, such as packet loss, jitter, and delay. An operator using the external probe may view results of a trace to diagnose the network communications problem. These results are not accessible to the network communications devices and cannot be accessed by network communications devices to alter network communications.\nTelecommunications switching systems today provide for Internet protocol (IP) communications between two end-points within a network or a different network to be terminated to a far end-point. Calls between two end-points are routed to the terminating end-point based on the address input at the originator. This address information is then relayed to a Call Control Manger (CCM) that screens, translates, and routes the call to the terminating subscriber or to another network to be terminated at a far end subscriber's end-point. The basic functionality of this process is widely known within the art and is used throughout telecommunications networks for voice calling.\nWithin the architecture of this switching system, calls to and from end-points are controlled by the CCM. The CCM may be located within a monolithic device in a TDM switch architecture or provided by an outboard computing device that controls the calls by using signaling that controls network based routing and switching devices located within the network. The latter device is known as soft-switch architecture.\nThe soft-switch architecture within an IP network controls call processing through use of signaling to and from the end devices and media gateways. One example of a protocol used for this IP signaling is Session Initiation Protocol (SIP). This protocol is currently used mainly with IP telephony, such as VoIP, and can be used as an access protocol between the end-user and the CCM and/or between the CCM of one network and the CCM of another network.\nAnother protocol used mainly between the CCM and a media gateway is the ITU-T H.248 protocol, commonly known as Megaco. This protocol is a control protocol that allows the CCM to control ingress and egress from/to the media gateway as calls are set up using a media gateway. Within a packet network framework, IP communications between two end-points (both access end-points and media gateways) are controlled by the signaling of the end-point to/from the CCM. The CCM provides authentication, screening, translations and routing based on information that is stored in the CCM and from the state of the end-points that the CCM controls.\nWithin the soft-switch architecture, call control can only be accomplished based on information possessed by the CCM or the on/off state of the devices that has an association with the CCM. While this configuration is fine in a static environment, packet networks are in a state of change at all times since the network itself can carry different types of information besides voice calls. One skilled in the art knows that a packet network is a converged network that can carry voice, data, and video all in a single path, and routing of calls within a packet network is not static and can vary significantly from call to call.\nBecause of packet network content communications variables, calls may encounter congestion and loss of voice quality based on latency, jitter, and packet loss. These content communications variables can affect any portion of a call at anytime based on the network elements usage at the time of the call. Unlike a TDM system where dedicated channels and circuits are provided, the CCM only has control of it own end-points. Other end units may attempt calls, computers may send/receive data without talking to the CCM, and other devices may require bandwidth while the original call is progressing, thus causing voice quality problems for the participants. In addition to these basic gaps, many physical layer 1 systems that are poor in regulating bandwidth, are being used for transmission facilities. WiFi, EVDO, 4G (Wax), DSL, and cable systems are all physical layer 1 technologies that demonstrate different bandwidth rates and management of their ability to modify available bandwidth as the Signal-to-Noise (SNR) ratios change.\nConventional soft-switching has not been designed to provide relief for callers when congestion, jitter or delay problems, such as those described above, are encountered. Since conventional CCMs can only determine call success based on connectivity to and from the calling parties, voice quality between two parties is not taken into consideration for call success.\nCommunication problems of in-band signals over packet networks are difficult to isolate. Currently, if a communication problem exists over a transmission path, there are few techniques to isolate the problem. One technique includes using an external probe to capture and decode packets, commonly know as a trace, traversing over a communications path to help isolate the problem. However, technicians generally only run the trace in response to a customer notifying a communications carrier of a communications problem. If a problem exists across packet networks operated by different carriers, a typical response by a carrier is to contact the other to determine if the other carrier can locate a problem in its network. In other words, locating an in-band communications problem over one or more packet networks is difficult as troubleshooting tools for such problems are limited to out-of-band performance metrics (PM) and are not available as in-band information via control or signaling paths.\nA problem that exists with current implementations of telephony over packet networks is that a call control manager does not have information about the bearer path. Traditionally, there was a linkage between transmission path state and the monolithic switch that essentially owned one end of that path where the in-band signaling and line characteristics were available and was an integral part of the information used by the CCM for call processing. As demonstrated in current implementations of VoIP, without knowledge of the bearer path, the call control manager may establish calls that result in poor voice quality or call setup failure.\nIn addition, IP Service gateways, such as a broadband remote access server (BRAS), functions to limit, commonly known as traffic rate shaping, each customer's DSL traffic to their purchased speed. There is no end-to-end signaling, outside of the embedded TCP flow control mechanism, used to adjust the bursting to eliminate packet loss. Rate shaping is a statically forced bandwidth constraint that alters the nature of a transmission path in the packet networks. This shaping coupled with commonly shared or “over-subscribed” bandwidth normally associated with trunking facilities between networks results in unknown transmission path states between media gateways servicing VoIP and other real-time services, such as Video on Demand (VOD).\nTraffic Quality of Service (QoS) management of packets is performed, where multiple flows aggregate into a smaller flows or channels. The application of Internet Protocol QoS is performed at the egress point where traffic is transmitted over a single link. Current traffic engines use the following information to make QOS traffic decisions. The decisions are assigning a Class of Service (CoS) and then acting upon that service to shape, restrict, or pass traffic to an egress point. The variables used to assign priority to traffic flows can be based on: entrance port (assign a whole port a CoS), virtual circuit in a port (assign a CoS to an Ethernet Virtual Circuit, LSP, etc.), priority bit marking of each packet (P bit), protocol type (assigning a CoS to specific types of packets or traffic), IP address and port (assigning a CoS to a whole IP address, or its port addresses), session identification (a HTTP, UDP, or other session addressed call), or otherwise. This priority marking information is used by service points, and shared links to implement QoS for the shared traffic flows. QoS and CoS types of information are made available at the point of aggregation where traffic management or QoS functions occur. However, the number of packets transmitted or lost in the packet stream elsewhere in the network is currently unavailable without the use of a session or path based protocol. These packet loss functions are generally not tracked by QoS mechanisms.\nIn current traffic rate shaping designs, the Internet may burst a packet stream to a DSL user when the packet network or Digital Subscriber Line Access Multiplexer (DSLAM) itself may not have sufficient bandwidth to accommodate the packet session. In a TCP-based session, the transmission rate is throttled down after packet loss is detected in the session. In VoIP, the packet loss is not counted by the use of Real Time Control Protocol (RTCP) signaling, but it is captured as the call progresses by the end points. RTCP, however, only considers performance of its own sessions and not the transmission path performance as a whole. In both cases, packets are sent over the packet network that get dropped in mid-path and will not make it to the customer premises equipment (CPE) and user. More importantly, there is no cross-session information about the packet loss and no whole path information available in-band.\nAlso, packet loss can be due to available bandwidth transmission rate fall-off, such as when a WiFi user walks away from a WiFi Access Point (AP) and loses RF signal strength, signal-to-noise ratio increases, or congestion increases due to many users concurrently accessing the AP. All of these types of conditions in the transmission paths can have severe impacts upon the ability to accomplish call processing and call management.\nAn Internet Service Provider (ISP) may provide different Internet connectivity speeds or data transfer rates based upon their service plans. For example, a user may purchase 1.5 Mbits/sec data transfer rate for a predetermined amount, such as 10 Mbits/sec data transfer rate or higher. In general, the transmission path is between the shared (trunked) BRAS resource and the DSLAM that is supplying. The normal amount of bandwidth consumption in the download direction from the network to the user is high as compared to the upstream direction. However, there is no correlated throttling mechanism in the IP web-server linked to user's ISP service plan that can be used to shape the packet transmissions. So, all of the network traffic is shaped at the BRAS typically based on the user's purchased data transfer rates. Depending upon network conditions, some of this traffic may not make it to the DSL user since the BRAS has no knowledge of the IP service path from itself to the customer.\nA problem occurs when the BRAS does realize congestion on a packet network where packets are being dropped due to insufficient bandwidth. Some packets could be dropped in the packet network or at an aggregation device somewhere in the packet network. Currently, there is little intelligence that recognizes the dropped packets in the packet network. In fact, packet networks are designed to discard traffic based on QoS markings. This problem is made worse because transporting packets that will ultimately be dropped adds to congesting the network. The packets consume bandwidth until dropped.\nTransmission Control Protocol (TCP) was designed to work in a best-effort, packet store-and-forward environment characterized by the possibility of packet loss, packet disordering, and packet duplication. Packet loss can occur, for example, by a congested network element discarding a packet. Here, a microprocessor or memory of a network element may not have adequate capacity to address all packets routing into and out of the element. Packet disordering can occur, for example, by routing changes occurring during a transmission. Here, packets of TCP connection may be being arbitrarily transmitted partially over a low bandwidth terrestrial path and as routing table changes occur partially over a high bandwidth satellite path. Packet duplication can occur, for example, when two directly-connected network elements use a reliable link protocol and the link goes down after the receiver correctly receives a packet but before the transmitter receives an acknowledgement for the packet.\nAn embedded capability within TCP protocol is the TCP sliding window technique. The sliding window was developed and deployed as a flow control mechanism used to minimize the inefficiencies of packet-by-packet transmission. The sending of data between TCP enabled end-devices on a connection is accomplished using the sliding window technique. TCP requires that all transmitted data be acknowledged by the receiving host. The sliding windows method is a process by which multiple packets of data can be affirmed with a single acknowledgement."}
{"text": "Modern application-specific integrated circuits (ASICs) integrate greater and greater security and data protection functionality into the hardware (HW). The integrated functionality provides more reliable and more efficient hardware security for both conventional “Data At-Rest” and conventional “Data In-Flight” protection.\nData storage systems are moving to distributed storage models that are based on storage networking. The move has an impact for enterprise data protection: the distributed models increase the vulnerability of stored data (i.e., Data At-Rest) to various attacks, both external and internal and both malicious and accidental. For Internet traffic and other moving data (i.e., Data In-Flight), the move provides such protection as sender and recipient mutual authentication, key exchange, data confidentiality, authenticated encryption (which is a type of encryption/decryption that additionally providing a way to check data integrity and authenticity) and replay protection.\nIn contemporary applications, the speed/throughput of the traversing data is up to 10 Gb/s (gigabits per second) and beyond. For some storage applications, the speed/throughput of the traversing data is even 10× higher: up to 100 Gb/s and beyond. The high speeds alone make security support of the data in software (SW) almost infeasible as far as security transformations are usually incorporated into the main data path and appear as bottlenecks from efficiency and performance standpoints.\nMany cryptographic protocols use an encryption process and message authentication and data integrity services independently with each process using an independent key. To speed up overall computations, new cryptographic modes that combine and provide both crypto services using a single “combined” mode were proposed and became accepted by both the National Institute of Standards and Technology (NIST) and the Institute of Electrical and Electronics Engineering (IEEE) and other technical professional organizations and committees working in network and data storage security areas.\nTo prevent data lost and breach, IEEE P1619 “Standard Architecture for Encrypted Shared Storage Media” suggests using the XTS-AES (Advanced Encryption Standard) (XOR-Encrypt-XOR (XEX)-based Tweaked Electronic Code Book (ECB) mode with Cipher Text Stealing (CTS)). The P1619.1 “Standard for Authenticated Encryption with Length Expansion for Storage Devices” uses the Galois/Counter mode (GCM), Counter mode (CTR) with Cipher-Block Chaining (CBC)-Message Authentication Code (MAC) (CCM) and other cryptographic processes. Both drafts are now accepted standards: IEEE Std. 1619-2007 and IEEE Std. 1619.1-2007.\nAmong the new AES-based modes is the NIST approved (see NIST Special Publication SP800-38D defining Galois/Counter Mode (GCM) and Galois Message Authentication Code (GMAC)) GCM mode and IEEE P1619 legacy mode Liskov, Rivest, and Wagner (LRW), that both use Galois Field multiplication for processing 128-bit blocks of data. Besides memory and storage applications, GCM-AES is becoming more widely used in various Internet security protocols and was suggested/submitted as an Internet-draft to the Internet Engineering Task Force (IEFT) to use in the Secure RTP (SRTP) protocol (see Internet-Draft for GCM in Secure RTP (SRTP)), MACsec (see IEEE 802.1AE), Internet Key exchange version 2 (IKEv2), and in the IPsec (see RFC 4106 and RFC 4543).\nA feature of the GCM mode is that the message authentication is performed in parallel with encryption/decryption of the main data payload by applying multiplication in a Galois Field (GF). Multiplications in finite fields have been used for fast (i.e., insecure) message hash computations. To make such computed massage hash values secure, application of the GCM GHASH process adds a pseudorandom vector, a so called “whitening” vector, at the end. The pseudorandom vector is generated by encrypting a preset value (i.e., Initialization Vector IV) with a secret AES key (i.e., vector W). Use of the GF multiplier for Message Authentication Code (MAC) computation permits higher throughput than the authentication process for computing a conventional MAC. The conventional MAC processes use slower chaining modes, like AES-CBC, or use a separate stand-alone secure hash process from the Secure Hash Algorithm (SHA) family."}
{"text": "The present invention relates generally to retention of one or more expansion cards installed within a computer chassis.\nPeripheral Component Interface (PCI) cards are often used to add functions and capabilities to computers. These cards include card edge contacts and are installed in a computer by insertion into a mating slot on a circuit board of the computer. These cards may extend perpendicularly from the circuit board and some means of support the card may be provided by an adjacent portion of the computer chassis. The PCI cards may be releasably attached to the means of support to prevent the cards from being disturbed from the slot and to prevent adjacent cards from contacting each other. Often, the releasable attachment is accomplished by the use of tools and typically some sort of removable fastener, such as a screw. A bracket may be attached to one end of the card and screwed to a card support or a wall of the chassis. Improvements to the attachment and release of expansion cards from slots within a computer chassis are desirable.\nPCI cards and the computers which receive them have been configured to permit hot swapping or plugging. This allows a card to be removed or installed from a slot in the computer without shutting the computer down. To accomplish this hot plug, the power supplied to the slot must be turned off prior to the removal or insertion of a card. Improvements to the powering or depowering of card slots for the insertion or removal of cards from the slots are desirable.\nEach card slot for receiving a PCI card within a computer chassis may have a first LED indicating whether the slot is currently powered up, and a second LED indicating the fault status of the slot. Often, these LEDs are mounted to the circuit board to which the slot is mounted and are located adjacent the card slot. Space on the circuit board may be crowded adjacent the card slot, making relocation of the LEDs desirable. Further, LEDs placed on the circuit board may be difficult to see from outside the chassis. It is desirable to improve the visibility of these LEDs while keeping them adjacent to the card slot they are associated with.\nThe present invention relates to a card retention device which includes a frame, a slider and a flap. The frame defines an opening through which an expansion card may be inserted into or removed from an expansion slot. The flap is hingedly mounted to the frame and movable between a captive position where the flap covers the opening and a free position where the flap does not cover the opening. The slider is slidably mounted to the frame and movable between a locked position where the slider holds the flap in the captive position and an unlocked position where the flap is free to move from the captive to the free position.\nThe present invention further relates to a computer chassis including an expansion slot for receiving an expansion card and a card support extending adjacent the slot. A card retention device is mounted to card support and includes a frame, a slider and a flap. The frame defines an opening through which an expansion card may be inserted into or removed from an expansion slot. The flap is hingedly mounted to the frame and movable between a captive position where the flap covers the opening and a free position where the flap does not cover the opening. The slider is slidably mounted to the frame and movable between a locked where the slider holds the flap in the captive position and an unlocked position where the flap is free to move from the captive to the free position."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital camera having autofocus capability and continuous shooting capability.\n2. Description of the Related Art\nA continuous shooting mode in which a series of photos are taken while the release button is held fully depressed (i.e., while the release switch is held ON) is known in the art as a shooting mode of a digital camera (see Japanese Unexamined Patent Publication 2007-017787). In conventional digital cameras, the camera is in an AF lock state (in which autofocus (AF) process is suspended) during continuous shooting. Therefore, even if the camera is in an in-focus state upon the commencement of continuous shooting, the camera sometimes becomes out of focus in the middle of continuous shooting when, e.g., shooting a moving object.\nTo keep the camera in focus during continuous shooting, an AF process only needs to be performed each time a picture is taken, i.e., each time the main mirror (quick-return mirror) moves down to the photographing position. However, a problem exists with this control slowing down the continuous shooting speed by the amount of time required to perform the AF process."}
{"text": "An auto document feeder (ADF) is a paper feeding mechanism for a photocopying machine. The ADF can convey a whole stack of paper sheets one by one to make them pass through a scanner of the photocopying machine for scanning operations. A common photocopying machine can only scan one side of the paper sheet. By adding another scanner in the ADF, the other side of the paper sheet can also be scanned, so both sides of the paper sheet can be scanned in one paper feeding operation.\nA simple configuration method, like disclosed in U.S. Pat. No. 7,457,006, U.S. Pat. No. 5,280,368, U.S. Pat. No. 7,411,704 and U.S. Pat. No. 8,730,537, is to include another scanner adjacent to the scanner of the photocopying machine. Generally, the scanner is a contact image sensor (CIS) scanner, so the paper sheet has to be in flat contact with the scanner for precise focus operations and better scanning quality. The two scanners are usually disposed opposite to each other. Therefore, after passing through one scanner, the paper sheet has to turn over instantly to come into flat contact with the other scanner. This motion can be achieved by utilizing a block for pressing paper sheets. However, this tends to cause twisting of the paper sheets and paper jams.\nAs a result, the target of the inventor is to solve the above-mentioned problems, on the basis of which the present invention is accomplished."}
{"text": "In general, a communication conference system is configured such that audio of conference participants is collected by a plurality of microphones, and is mixed and then transmitted. On the reception side, the audio is reproduced from all of loudspeakers at an equal level of sound volume and phase. In the case of such a communication conference system, the audio data is reproduced with a same sound image regardless of who is speaking among the conference participants. As a result, it is difficult on the reception side to identify who spoke.\nMoreover, on the reception side, since the sound volume of a loudspeaker is fixed, the sound volume from the loudspeaker significantly fluctuates when there is a plurality of conference participants on the transmission side, due to differences in the audio sound volume of the respective participants and in the distance between microphones and the participants.\nBased on such circumstances, there has been proposed a voice telephone conference device that determines a speaker and the position of the speaker based on the temporal waveform and the frequency spectrum of a microphone input signal (for example, refer to Patent Document 1).                Patent Document 1: Japanese Unexamined Patent Application, First Publication (JP-A) No. H09-261351        \nThe voice telephone conference device disclosed in Patent Document 1 determines a speaker and the position of the speaker based on the temporal waveform and the frequency spectrum of a microphone input signal, and transmits position information along with audio data. The reception side controls the sound volume of each loudspeaker based on the received position information. Moreover, there is provided a switch for switching to select audio data with a highest input level among the respective microphones, and the reception side is set so that the audio is emitted from the loudspeaker that corresponds to the microphone that collected the audio.\nHowever, a sound field with a high level of presence (for example, a feel of depth) can not be realized by only controlling the sound volumes of the respective loudspeaker (stereo loudspeaker units) based on the position information.\nFurthermore, each of the microphones is provided with the switch for switching to select audio data with a highest input level among the respective microphones. However, there is an issue in that when a plurality of speakers start speaking simultaneously, this switching needs to be performed in a short period of time, and accurate audio data cannot be transmitted as a result."}
{"text": "1. Field of the Invention\nThe present invention relates to a retractor for use in a vehicle seat belt system and particularly, to a seat belt retractor adapted to inhibit extraction of a seat belt in emergency situations such as collision.\n2. Description of the Related Art\nA motor vehicle conventionally includes a seat belt system mounted to a passenger's seat to protect a vehicle occupant in emergency situations such as collision. In such a seat belt system, an emergency lock type retractor is attached to a rigid member to take up a seat belt.\nSuch a retractor is designed to allow extraction of a seat belt when the seat belt is fastened around a vehicle occupant. The seat belt can also be extracted in non-emergency situations so as not to restrain movement of the occupant. In emergency situations such as collision, impact or sudden acceleration is sensed to actuate a reel lock mechanism so as to lock a reel around which the seatbelt is being wound. This allows the seat belt to restrain the vehicle occupant so as to inhibit sudden movement of the occupant or to protect the occupant.\nThe reel lock mechanism is operable to securely stop the reel per se around which the seat belt is wound. However, the seat belt may be extracted in the event that it is loosely wound around the reel. To prevent this, there has previously been proposed a seat belt retractor as shown in FIGS. 3 to 6.\nAs shown in FIGS. 3 and 4, a frame 10 includes a pair of parallel side walls 12 and 14, and a rear wall 16 extending between the side walls 12 and 14. A reel 20 and an emergency lock mechanism 22 are mounted to the lower portion of the frame 10 to take up a seat belt 18. A belt lock mechanism 24 is mounted to the upper portion of the frame 10 to inhibit extraction of the seat belt 18 at the time of emergency.\nThe reel lock mechanism 22 will now be described with reference mainly to FIG. 4.\nThe side walls 12 and 14 include coaxial support openings 26 and 28 through which a reel shaft 30 extends while passing through a bushing 32 made of synthetic resin and is rotated about its own axis.\nA return spring 34 is connected at a central end to one end of the shaft 30. A cover 36 is disposed over the return spring 34 to secure one side of the return spring 34 to the side wall 12. Extraction of the seat belt 18 causes the reel 20 to rotate in a direction to store energy in the return spring 34. When the occupant releases his hand from the seat belt 18, the reel 20 is rotated under the biasing force of the return spring 34. Then, the seat belt 18 is automatically wound around the reel 20. It will be understood that the reel 20 and the reel shaft 30 are rotated in the direction of the arrow A.sub.1 when the seat belt 18 is extracted.\nThe reel lock mechanism 22 is mounted to the outer surface of the side wall 14.\nThe reel lock mechanism 22 includes a ratchet wheel 38 integral with the other end of the reel shaft 30. The ratchet wheel 38 includes a pin 40 coaxial with the shaft 30. A tie plate 42 has an opening 43 fitting around the pin 40. A lock ring 44 has a central opening 48 in which the pin 40 is loosely fit. An arcuate spring element 50 has one end engaged with a central hole (spring hanger) 52 of the tie plate 42 and the other end engaged with a hole (spring hanger) 54 of the lock ring 44. The lock ring 44 has internal teeth 56. The spring element 50 extends between the spring hanger 54 of the lock ring 44 and the spring hanger 52 of the tie plate 42 and provides a biasing force to rotate the lock ring 44 in the direction of the arrow A.sub.2.\nA control lever 58 has a base end pivotally connected to the side wall 14 of the frame 10 by a pivot pin 60. The other, free end of the control lever 58 is engageable with the ratchet wheel 38. A pin 62 extends from one side of the control lever 58. The pivot pin 60 extends through an opening 64 which is formed in the front end of the tie plate 42.\nThe lock ring 44 has a pair of diametrically opposite integral tabs 66 and 68. The tab 66 is designed to rotate the control lever 58, whereas the tab 68 is designed to operate the belt lock mechanism 24.\nThe tab 66 of the lock ring 44 has an elongate hole 70 to receive the pin 62 of the control lever 58.\nA hook retainer 72 is secured to the pin 40 of the shaft 30 which in turn, extends through the central opening 48 of the lock ring 44. A pair of diametrically opposite projections 76 and 78 extend from the peripheral edge of the hook retainer 72 to support a hook 74. The hook 74 has two openings 80 and 82 to receive the projections 76 and 78. This arrangement allows the hook 74 to reciprocate on a line extending between the projections 76 and 78 (shown by the arrows B.sub.1 and B.sub.2).\nA compression coil spring 84 is disposed between the hook retainer 72 and the hook 74 to urge the hook 74 in the direction of the arrow B.sub.1. A pawl 86 extends from the outer peripheral edge of the hook 74 to engage with the internal teeth 56 of the lock ring 44. A connecting pin 88 extends from one side of the hook 74.\nThe hook 74 is normally urged in the direction of the arrow B.sub.1 by the compression coil spring 84. That is, the hook 74 is shifted to the left as shown in FIG. 4. This results in separation of the pawl 86 from the internal teeth 56.\nA frictional engagement member 90 is substantially in the form of a ring. An opening 92 is formed adjacent to the outer peripheral edge of the frictional engagement member 90 to receive the connecting pin 88. A flywheel 96 is fit over the frictional engagement member 90 and includes a ratchet 94. The flywheel 96 has a central opening within which the pin 40 of the shaft 30 is loosely fit. The flywheel 96 is short and cylindrical in shape. The frictional engagement member 90 is fit within the flywheel 96. An arcuate spring 98 is fit on the outer periphery of the frictional engagement member 90 and urged against the inner periphery of the flywheel 96. Friction between the flywheel 96 and the frictional engagement member 90 enable sliding rotary motion of the flywheel 96. A hole 40A is formed in the leading end of the pin 40 to receive a rivet 96. This holds the flywheel 99 in place.\nAs shown in FIG. 3, an actuator 100 is mounted to the side wall 14 of the frame 10 and generally includes a case 104 fixed to the side wall 14, an operating element or barrel 106 loosely received in the case 104, an operating piece 110 having a protrusion 108 in contact with the upper surface of the operating barrel 106, and a support 112 by which the base end of the operating piece 110 is pivotally supported.\nAs shown in FIG. 3, a cover 114 surrounds the reel lock mechanism assembled in a manner shown in FIG. 4.\nWith the retractor of the seat belt thus constructed, the operating piece 110 is disengaged from the flywheel 96 when the seat belt is extracted by the vehicle occupant. This permits rotation of the reel 20 and the shaft 30 and thus, extraction of the seat belt 18. If the seat belt 18 is released, then the shaft 30 is rotated under the influence of the return spring 34 within the cover 36 so that the seat belt 18 may be wound around the reel 20.\nIf the speed of the vehicle is substantially changed due, for example, to collision, then the actuator 100 is rendered operative to inhibit extraction of the seat belt 18. Specifically, the operating barrel 106 is inclined when a substantial amount of acceleration is exerted on the actuator 100. This causes the protrusion 108 to push up the operating piece 110. A free end of the operating piece 110 is then brought into engagement with the ratchet 94. As a result, the flywheel 96 is prevented from rotating.\nStoppage of the flywheel 96 results in corresponding stoppage of the frictional engagement member 90. A vehicle collision causes extraction of the seat belt 18. This would result in rotation of the reel shaft 30 and thus, the hook retainer 72 and the hook 74 rotate in the direction of the arrow A.sub.1. However, the hook 74 can not be rotated since the frictional engagement member 90 is prevented from rotating as a result of engagement with the pin 88. The hook 74 is slid in the direction of the arrow B.sub.2 to the extent corresponding to the rotation of the hook retainer 72 in the direction of the arrow A.sub.1. The pawl 86 is then brought into engagement with the internal teeth 56 of the lock ring 44.\nConsequently, the lock ring 44 is rotated in the direction of the arrow A.sub.1. Rotation of the tab 66 in the direction of the arrow A.sub.1 causes the control lever 58 to rotate in the direction of the arrow C.sub.1 since the pin 62 is engaged with the elongate hole 70. The free end of the control lever 58 is then brought into engagement with the ratchet wheel 38 of the reel shaft 30 so as to firmly lock the reel shaft 30 and the reel 20.\nThe construction of the belt lock mechanism 24 will next be described with reference to FIGS. 3, 4 and 5.\nA first gripping member 120 is attached to the rear wall 16 of the frame 10 and includes a holder 123 vertically movable along the rear wall 16, a receiver 124 held by the holder 123, and a spring 126 disposed to urge the receiver in a downward direction. The receiver 124 has a rugged front surface.\n122 is a guide member for guiding the seat belt 18. As shown in FIG. 6, the guide member or frame 122 has upper and lower slots 127 and 129 and is made of synthetic resin.\nA pair of openings 130 and 132 are coaxially formed in the side walls 12 and 14 of the frame 10 to receive the shaft 134. A second gripping member 136 has a hole 137 through which the shaft 134 extends. Thus, the second gripping member 136 is rotatably mounted to the frame 10.\nAs shown in FIG. 5, a semicylindrical pusher 138 is attached to the free end of the second gripping member 136 and has a rugged surface in a face-to-face relation to the receiver 124. The seat belt 18 extends between the pusher 138 and the receiver 124.\nA pin 140 (FIG. 3) extends from one side of the second gripping member 136 into an elongate hole 142 which is formed in the side wall 14 of the frame 10.\nA rocker arm 146 is pivotally mounted to the outer side of the wall 14 of the frame 10 by a pivot pin 144. The rocker arm 146 is L-shaped and has a notch 148 at one end to receive the pin 140. The other end of the rocker arm 146 is pivotally connected to the upper end of a lever 150 to form a joint 152. The lever 150 is connected to the joint 152 of the rocker arm 146 in a manner to allow slight angular movement in a direction as indicated by E.\nThe lower end of the lever 150 is overlapped with the tab 68 of the lock ring 44. A pin 154 extends from the lower end of the lever 150 into an elongate hole 156 of the tab 68.\nA spring 158 is mounted to the pivot pin 144 and has one end engaged with a hole 160 of the side wall 14 of the frame 10 and the other end engaged with a hole 162 of the rocker arm 146. The spring 158 urges the rocker arm 146 in the direction of the arrow G.sub.1.\nOperation of the belt lock mechanism 24 thus constructed is as follows.\nIn a non-emergency situation, the rocker arm 146 is urged in the direction of the arrow G.sub.1 under the action of the spring 158 to push the pin 140 in the same direction. This causes the pusher 138 of the second gripping member 136 to separate from the receiver 124 so as to allow passing of the seat belt 18 between the pusher 138 and the receiver 124.\nIn emergency situations such as collision, the operating barrel 106 of the actuator 100 is inclined to cause the operating piece 110 to engage with the ratchet 94 of the flywheel 96 as mentioned earlier. As a result, the lock ring 44 is rotated in the direction of the arrow A.sub.1. Then, the seat belt 18 will be locked in the following steps 1 to 5.\n1 Rotation of the tab 68 of the lock ring in the direction of the arrow A.sub.1 results in downward movement of the pin 154 within the elongate hole 156, and thus, rotation of the lever 150 in a direction as indicated by E.\n2 The rotation of the lever 150 causes the free end 150a to engage with the ratchet wheel 38.\n3 When the ratchet wheel 38 is rotated, the lever 150 is moved upwards. The rocker arm 146 is then rotated about the pivot pin 144 in a direction as indicated by G.sub.2. The pin 140 is pushed in a direction as indicated by G.sub.2 (The elongate hole 70 is so shaped that when the pin 140 is pushed in such a direction, the control lever 58 must not engage with the ratchet wheel).\n4 As a result, the second gripping member 136 is moved toward the first gripping member 120 so as to grip or sandwich the seat belt 18 between the pusher 138 and the receiver 124.\n5 Once the seat belt 18 is sandwiched between the pusher 138 and the receiver 124, extraction of the seat belt 18 urges the pusher 138 and the receiver 124 toward one another. As a result, the seat belt 18 is firmly gripped between the pusher 138 and the receiver 124 and can no longer be extracted.\nIn the prior art seat belt retractor, the seat belt 18 tends to slightly slide between the two gripping members 120 and 136 when the first gripping member 120 and the second gripping member 136 begin to grip the seat belt 18.\nTo present sliding motion of the seat belt 18, the pusher 138 or the receiver 124 could have a serrated surface to contact the seat belt. However, such a serrated surface presents a problem that the seat belt is rubbed to create a rough nap."}
{"text": "Hoses are often reinforced with continuous yarn to improve physical performance characteristics, such as burst strength. For example, radiator hoses for automobiles and trucks are generally reinforced with continuous yarn reinforcing elements. Continuous yarn reinforcements are typically used to improve the burst strength of such hoses. Even though building such fiber-reinforced hoses is a labor-intensive operation which results in a substantial amount of material waste, such techniques have been required to meet the demands of the automotive industry.\nBuilding fiber reinforcements into hoses is a labor-intensive operation. After such hoses are built, they are typically trimmed to the exact size required. The fiber reinforcement containing material trimmed from such hoses generally has to be scrapped since it is not typically possible to recycle such fabric containing material.\nHoses have been made by extruding rubber compositions into the form of a tube which is subsequently shaped into the desired form and cured. Such techniques are advantageous in that they reduce labor costs, in-process inventory and waste. However, hoses made utilizing such extrusion techniques have typically not had the physical strength demanded by the automotive industry for radiator hose. More specifically, such hoses have had low burst strength.\nBy utilizing the technique disclosed in U.S. Pat. No. 5,268,134, hose which exhibits a burst strength of about 115 pounds per square inch can be made without utilizing fabric reinforcements. This process specifically comprises (1) extruding a rubber composition into the form of a tube, wherein the rubber composition is comprised of (a) an EPDM alloy comprised of (i) functionalized EPDM rubber, (ii) an EPDM rubber having thermoplastic side chains grafted thereto and (iii) dispersed thermoplastic wherein the thermoplastic is selected from the group consisting of nylons, polyesters and polyphenylene oxides, (b) an EPDM rubber, (c) carbon black, (d) at least one curative, (e) zinc oxide, (f) a processing oil and (g) stearic acid; (2) shaping the tube into the geometric form desired for the hose; and (3) curing the rubber composition at a temperature within the range of 130.degree. C. to 210.degree. C. to produce the hose.\nU.S. Pat. No. 5,268,134 also discloses a process for manufacturing hose which comprises injection molding a rubber composition into the desired geometric form for the hose at a temperature which is within the range of 130.degree. C. to 210.degree. C., wherein the rubber composition is comprised of (a) an EPDM alloy comprised of (i) functionalized EPDM rubber, (ii) an EPDM rubber having thermoplastic side chains grafted thereto, and (iii) dispersed thermoplastic wherein the thermoplastic is selected from the group consisting of nylons, polyesters, and polyphenylene oxides, (b) an EPDM rubber, (c) carbon black, (d) at least one curative, (e) zinc oxide, (f) a processing oil, and (g) stearic acid; (2) shaping the tube into the geometric form desired for the hose; and (3) curing the rubber composition at a temperature within the range of 130.degree. C. to 210.degree. C. to produce the hose."}
{"text": "As part of the rapid growth of Internet and World Wide Web use, there has been an ever increasing growth in the availability of online services. Such online services include, for example, online banking and financial services, online email services, online retail services, online dating services and online social networks. A given online service provider may provide a number of such services, which users of the online services typically access using a username and password, or other user login credentials that may be used by a provider as a mechanism for verifying users of the provider's services.\nAccessing online services using such login credentials is, however, subject to fraud and abuse. For example, malicious actors may obtain user login credentials using improper means. Such improper means may include the use of malicious software, such as could be installed on publicly available computers (e.g., in libraries, Internet cafes or other locations), to monitor keystrokes on an affected computer to capture users' login credentials. User credentials could also be improperly obtained using brute force attacks by trying various combinations of usernames and passwords through an automated trial and error process. These techniques for improperly obtaining login credentials are provided by way of example and a number of other approaches for improperly obtaining user credentials may be used. For example, login credentials may be improperly obtained from a user through a process referred to as phishing. Phishing may be accomplished by malicious actors that masquerade as a trustworthy entity (e.g., a bank or credit card company) by sending, for example, a fraudulent email message that appears to be from a trustworthy entity. These phishing messages may include a request that deceives a user into providing login credential information to a malicious actor.\nOnce a malicious actor has acquired a user's login credentials for a given online service provider, the malicious actor may then gain access to the user's account with that provider. Depending on the specific provider and online services provided by that provider, the malicious actor could then make fraudulent purchases using the user's financial information, gain access to the user's bank accounts, or perform other acts of identity fraud, posing as the legitimate user for whom the malicious actor has obtained login credentials."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for forming self-assembled colloidal photonic crystals on a selected area of a substrate or forming self-assembled colloidal photonic crystals having different diameters on the same substrate using an electric field, and a method for fabricating three-dimensional photonic crystal waveguides of an inverted-opal structure using the self-assembled colloidal photonic crystals patterned by the patterning method.\n2. Description of the Related Art\nPhotonic band gap structures in photonic crystals composed of dielectrics having a three-dimensional periodicity have become a matter of increasing interest. The photonic band gap structures are highly applicable to diverse optoelectronic devices such as microlasers, filters, high-efficiency LEDs, optical switches, low-loss waveguides, etc. In an initial stage of investigation, a three-dimensional photonic band gap in a microwave range was realized by providing periodicity to the dielectric constant of a dielectric-air structure by making small holes in a parallel direction on a silicon wafer or by stacking bar-shaped dielectrics in piles. However, in the wavelength ranges of infrared rays and visible light, it was only possible to form a two-dimensional photonic band gap. This is because it was very difficult to form a three-dimensional photonic band gap due to the need for scaling-down the etching space. In the case of a three-dimensional photonic band gap, the holographic lithography of laser light was used. Recently, research for methods of self-assembling small spheres (colloids) having a diameter of several hundred nanometers (nm) has been conducted.\nIn particular, diverse methods for fabricating self-assembled colloidal photonic crystals have been studied. One of the methods most frequently used is a dip-coating method for fabricating photonic crystals which utilizes capillary force exerted among a colloidal fluid, a substrate and colloidal particles. This dip-coating method is easy to carry out, and can form photonic crystals having a high crystallization in a wide area. However, it is difficult to selectively control the colloidal particles, and many semiconductor processes such as the lithography are required for patterning the photonic crystals. Additionally, in the case of forming photonic crystal of different kinds of colloidal particles or colloidal particles having different sizes, a template is required. Also, in the case of forming colloidal photonic crystals using three or more kinds of colloidal particles or colloidal particles having three or more sizes, the fabricating processes become complicated with limitations in design.\nThe most frequently used method for fabricating waveguides using photonic crystals comprises forming a transmission line from two-dimensional photonic crystals formed by periodically etching fine holes onto a silicon substrate. The two-dimensional photonic crystal has a photonic band gap and therefore results in no photonic loss in a two-dimensional light-traveling direction, but does suffer photonic loss in other light-traveling directions.\nMeanwhile, because the three-dimensional photonic crystal of an inverted-opal structure has a photonic band gap in all three-dimensional light-traveling directions, a waveguide using these three-dimensional photonic crystals can greatly reduce photonic loss in comparison to a two-dimensional photonic crystal. The fabrication of a waveguide having the three-dimensional photonic band gap requires a high-grade etching technique such as e-beam lithography, a high fabricating cost and considerable manufacturing time, and it is difficult to implement the waveguide in a wide area. Accordingly, there has been an increasing demand for and interest in a method for fabricating three-dimensional photonic crystal waveguides of an inverted-opal structure through self-assembly of colloidal particles that can easily allow for fabricating photonic crystals in a wide area.\nConventional methods for fabricating three-dimensional photonic crystal waveguides of an inverted-opal structure have been proposed as follows.\nFirst, the method described in “Multi-Photon Polymerization of Waveguide Structures within Three-Dimensional Photonic Crystals”, Advanced Materials, vol. 14, 2003, pp 271-294 by W. Lee et al. has problems due to material and area limitations. Second, the method described in “Micromolding of Three-Dimensional Photonic Crystals on Silicon Substrates”, Nanotechnology, vol. 14, 2003, pp 323-326 by P. Ferrand et al. also has problems in that its process is complicated and the photonic crystals are broken with the occurrence of cracks during an artificial piling up of the photonic crystals."}
{"text": "Field\nThe application relates to a medical apparatus, and more particularly, relates to a medical apparatus for delivering radiotherapy.\nDescription\nRadiotherapy is a method which transmits radioactive rays such as α-rays, β-rays, or γ-rays, generated by radioisotopes, X-rays, electron beams, proton beams or other particles, to diseased tissue (e.g., a cancerous tumor). The rays kill cells of the tissue by causing ionizations within the cells or other cell damage.\nA linear accelerator is a particle accelerator commonly used for radiotherapy. A linear accelerator comprises a radiation head in which a radiation source is arranged for radiating a radioactive beam toward diseased tissue. The radiation source may include an accelerating tube, an electron gun, a moving target, a magnetic biasing system, a collimator, and a flattening filter. The radiation head also includes a high-density lead shielding layer to prevent extraneous radioactive rays from harming bystanders during use of the linear accelerator. The above components, and others, contribute to the substantial mass of the radiation head.\nConventional radiotherapy is typically performed in conjunction with a computed tomography (hereinafter, “CT”) imaging apparatus or a magnetic resonance (hereinafter, “MR”) imaging apparatus to determine a specific position of the diseased tissue and to thereby properly position the diseased tissue relative to the emitted rays of the linear accelerator. As is known in the art, a CT imaging apparatus is particularly suited for imaging bone and an MR imaging apparatus is particularly suited for imaging soft tissue.\nIn one example, a patient is positioned within an imaging apparatus to image the diseased tissue and refine the patient's position, and the patient is then moved to the linear accelerator to radiate the diseased tissue. Since both the linear accelerator and the imaging apparatus occupy a large volume, the patient is moved over a substantial distance therebetween, thereby complicating the procedure and increasing the risk and extent of positioning errors.\nFIG. 19 illustrates a prior combined medical apparatus 100. The medical apparatus 100 comprises a radiation therapy assembly 102, an imaging assembly 104, and a couch assembly 106. More particularly, the radiation therapy assembly 102 is a common radiotherapy apparatus, which comprises a gantry 1022 in which a through-hole 1024 extending along a horizontal direction is defined. One end of an arm 1026 is secured to the gantry 1022 and the other end thereof extends outwardly. A radiation head 1028 is fixed to the other end of the arm 1026. The imaging assembly 104 is a common CT imaging apparatus, which comprises a gantry 1042 defining a through-hole 1044 extending along a horizontal direction. An X-ray tube 1046 and a detector 1048 are arranged oppositely on a rotatable mechanism (not shown) surrounding the through-hole 1044. The couch 106 comprises a base 1062 and a patient support 1064 arranged on the base 1062.\nIn order to image a patient, the patient is placed on the patient support 1064, and the patient support 1064 is moved through the through-holes 1024 and 1044 in order to position the patient between the X-ray tube 1046 and the detector 1048. The X-ray tube 1046 and the detector 1048 are then operated as is known in the art to acquire an image of the patient. To radiate the patient, the patient support 1064 is moved such that the target volume, determined from the acquired image, is positioned in the beam field of the radiation head.\nThe length spanned by the through-holes 1024 and 1044 or, more particularly, the distance from the isocenter of the radiation therapy assembly 102 to the imaging plane of the imaging assembly 104, is quite long. The resulting movable distance of the patient support 1064 is directly related to the likelihood that deformation of the couch assembly and/or the patient support will introduce errors in imaging and/or radiation delivery. Additionally, as recited above, the large volume occupied by the combined medical apparatus 100, consisting of the radiation therapy assembly 102 and the imaging assembly 104, poses a significant challenge to any institution (e.g., a hospital) employing the apparatus 100.\nFIG. 20 illustrates another prior combined medical apparatus 200. The medical apparatus 200 comprises a radiation therapy assembly 202, an imaging assembly 204 and a couch assembly 206. The radiation therapy assembly 202 comprises a gantry 2021 defining a through-hole 2022 extending along the horizontal direction. One end of an arm 2023 is secured to the gantry 2021 and the other end thereof extends outwardly. A radiation head 2024 is fixed on this other end of the arm 2023. The imaging assembly 204 is a CT imaging assembly which comprises a gantry 2041 defining a through-hole 2042 along the horizontal direction. An X-ray tube 2043 and a detector 2044 are oppositely arranged around through-hole 2042 on a rotatable mechanism (not shown). The couch assembly 206 comprises a base 2061 and a rotatable platform 2062 which is coupled to the base 2061 and rotatable about a vertical axis 2063. A patient support 2064 is slidably secured to the platform 2062, and a patient P is shown positioned on the patient support 2064.\nIn contrast to the combined medical apparatus 100 as shown in FIG. 19, the radiation therapy assembly 202 and the imaging assembly 204 of the medical apparatus 200 are separately arranged, and the couch assembly 206 is arranged between the radiation therapy assembly 202 and the imaging assembly 204. Accordingly, during the transition between an imaging mode and a radiating mode, the movement of the patient support 2064 differs from the movement of the patient support 1064 as shown in FIG. 19, i.e., the patient support 2064 rotates to change between the two modes but the patient support 1064 translates linearly to change modes.\nThe movable distance of the patient support 2064 is shorter than that of the patient support 1064 in FIG. 19, and thus, the apparatus 200 attempts to address the problem of flexible deformation of the patient support 2064. However, since the apparatus 200 requires sufficient distance between the radiation therapy assembly 202 and the imaging assembly 204 such that the patient support 2064 can be rotated without interference with the radiation therapy assembly 202 and the imaging assembly 204, the space occupied by the medical apparatus 200 is unsuitably large.\nThus, technical solutions are desired in radiotherapy to reduce complications, increase efficiency of resource usage, and decrease positioning error."}
{"text": "1. Field of the Invention\nThe present invention relates to alerting devices, and more particularly to a temperature based alerting device for obtaining a user\"\"s attention upon the occurrence of a predetermined event.\n2. Description of the Related Art\nThe use of wireless communication is increasing at an exponential rate, and as such there have been many attempts among the service providers and manufacturers of the hardware implemented in the wireless systems to provide options to the users that enable more convenient and more user friendly and discrete access to the services. Among these attempts have been the implementation of other methods for alerting the user of the occurrence of a particular event and more specifically, the occurrence of an incoming communication signal. For example, with radio paging devices, the user generally has two options for setting the alert type when an incoming page is received; (i) an audible tone to indicate the presence of an incoming page or alternatively, (ii) a vibration mode to vibrate the device when the incoming page is being received. The vibration mode of the radio pager serves the same function of obtaining the user\"\"s attention when an incoming page is received but does so in a silent manner so as to not disturb the user or persons surrounding the user when the page is received.\nSome wireless telephones also provide a similar vibration mode for indicating the presence of an incoming telephone call. The vibration mode has been implemented into these devices primarily to provide the user with the option of turning off the audible tones generated by the radio pager or wireless telephone to prevent inconvenient audible disturbances in a variety of different places. Examples of these places would be office meetings, libraries, and any other place or circumstance that the user deems appropriate to eliminate the audible tones generated by the respective devices.\nU.S. Pat. No. 5,861,686 to Lee discloses a device for generating waking vibrations or sounds. The device is implemented into alarm watches or in communications equipment such as cellular pagers or phones. The device utilizes an electromagnet, a element and a first and second set of vibration members in a ring case. These elements are connected to the printed circuit board PCB of the device (e.g. watch, phone or pager) and together enable the selective generation of vibratory motion when the device receives an incoming call, or in the case of a watch for providing an alarm function.\nU.S. Pat. No. 5,619,181 to Murray discloses a vibratory alerting device with audible sound generator. The alerting device simultaneously generates a vibration alert and an audible alert to notify the user as to the presence of an incoming call on a portable communication device such as a pager or wireless telephone.\nAs mentioned previously, the use of an audible alert signal can be undesirable when the user is located in place where an audible alert signal would be considered an interruption. In addition, the devices used to implement a vibratory alert signal in a communication device are generally bulky in nature and have mechanical moving parts which require additional space within the device and thereby increases the overall size of the device.\nFurthermore, all of the existing alerting devices (e.g., audible and vibratory) are contained within the communication device (i.e., phone or pager) and therefor require the user to be carrying the same in order to receive the alerting signal.\nIt would be advantageous to provide a more discrete alerting device that is not physically connected or disposed within the user\"\"s communication device. This would enable the user to be alerted as to the presence of an incoming communication signal without requiring them to carry the communication device.\nIn accordance with a preferred embodiment of a invention, the method for generating a temperature based signal for alerting a user to the occurrence of a predetermined event comprises: providing a temperature based alerting device having at least one heating element in contact with the user\"\"s skin; initializing the temperature based alerting device to a predetermined temperature; setting a temperature change limit for the temperature based alerting device; changing the temperature of the heating element of the temperature based alerting device upon occurrence of the predetermined event; and resetting the temperature of the heating element of the temperature based alerting device when the user has responded to the occurrence of the predetermined event.\nOther objects and features of the present invention will become apparent from the following detailed description considered in conjunction with the accompanying drawings. It is to be understood, however, that the drawings are designed solely for purposes of illustration and not as a definition of the limits of the invention, for which reference should be made to the appended claims. It should be further understood that the drawings are not necessarily drawn to scale and that, unless otherwise indicated, they are merely intended to conceptually illustrate the structures and procedures described herein."}
{"text": "Many companies and other organizations operate computer networks that interconnect numerous computing systems to support their operations, such as with the computing systems being co-located (e.g., as part of a local network) or instead located in multiple distinct geographical locations (e.g., connected via one or more private or public intermediate networks). For example, distributed systems housing significant numbers of interconnected computing systems have become commonplace. Such distributed systems may provide back-end services to servers that interact with clients. Such distributed systems may also include data centers that are operated by entities to provide computing resources to customers. Some data center operators provide network access, power, and secure installation facilities for hardware owned by various customers, while other data center operators provide “full service” facilities that also include hardware resources made available for use by their customers. As the scale and scope of distributed systems have increased, the tasks of provisioning, administering, and managing the resources have become increasingly complicated.\nThe advent of virtualization technologies for commodity hardware has provided benefits with respect to managing large-scale computing resources for many clients with diverse needs. For example, virtualization technologies may allow a single physical computing device to be shared among multiple users by providing each user with one or more virtual machines hosted by the single physical computing device. Each such virtual machine may be a software simulation acting as a distinct logical computing system that provides users with the illusion that they are the sole operators and administrators of a given hardware computing resource, while also providing application isolation and security among the various virtual machines. With virtualization, the single physical computing device can create, maintain, or delete virtual machines in a dynamic manner.\nWhile embodiments are described herein by way of example for several embodiments and illustrative drawings, those skilled in the art will recognize that embodiments are not limited to the embodiments or drawings described. It should be understood, that the drawings and detailed description thereto are not intended to limit embodiments to the particular form disclosed, but on the contrary, the intention is to cover all modifications, equivalents and alternatives falling within the spirit and scope as defined by the appended claims. The headings used herein are for organizational purposes only and are not meant to be used to limit the scope of the description or the claims. As used throughout this application, the word “may” is used in a permissive sense (i.e., meaning “having the potential to”), rather than the mandatory sense (i.e., meaning “must”). Similarly, the words “include,” “including,” and “includes” mean “including, but not limited to.”"}
{"text": "This present invention relates to a hydrocarbon conversion process. More particularly, this invention relates to the catalytic hydrocracking of hydrocarbons.\nThe hydrocracking of hydrocarbons is old and well-known in the prior art. These hydrocracking processes can be used to hydrocrack various hydrocarbon fractions such as reduced crudes, gas oils, heavy gas oils, topped crudes, shale oil, coal extract and tar extract wherein these fractions may or may not contain nitrogen compounds. Modern hydrocracking processes were developed primarily to process feeds having a high content of polycyclic aromatic compounds, which are relatively unreactive in catalytic cracking. The hydrocracking process is used to produce desirable products such as turbine fuel, diesel fuel, and middle distillate products such as naphtha and gasoline.\nThe hydrocracking process is generally carried out in any suitable reaction vessel under elevated temperatures and pressures in the presence of hydrogen and a hydrocracking catalyst so as to yield a product containing the desired distribution of hydrocarbon products.\nHydrocracking catalysts generally comprise a hydrogenation component on an acidic cracking support. More specifically, hydrocracking catalysts comprise a hydrogenation component selected from the group consisting of Group VIB metals and Group VIII metals of the Periodic Table of Elements, their oxides or sulfides. These hydrocracking catalysts contain an acidic support comprising a crystalline aluminosilicate material such as X-type and Y-type aluminosilicate materials. This crystalline aluminosilicate material is generally suspended in a refractory inorganic oxide such as silica, alumina, or silica-alumina.\nRegarding the hydrogenation component, the preferred Group VIB metals are tungsten and molybdenum; the preferred Group VIII metals are nickel and cobalt. The prior art has also taught that combinations of metals for the hydrogenation components, expressed as oxides and in the order of preference, are: NiO--WO.sub.3, NiO--MoO.sub.3, CoO--MoO.sub.3, and CoO--WO.sub.3. Other hydrogenation components broadly taught by the prior art include iron, ruthenium, rhodium, palladium, osmium, indium, platinum, chromium, vanadium, niobium, and tantalum.\nThere is a myriad of catalysts or catalyst systems known for hydrocracking whose properties vary widely. A catalyst suitable for maximizing naphtha yield may not be suitable for maximizing the yield of turbine fuel or distillate. Further, various reactions; i.e., denitrogenation, hydrogenation, and hydrocracking must be reconciled in a hydrocracking process in an optimum manner to achieve the desired results.\nFor instance when a feedstock having a high nitrogen content is exposed to a hydrocracking catalyst containing a high amount of cracking component the nitrogen serves to poison or deactivate the cracking component. Thus, hydrodenitrogenation catalysts do not possess a high cracking activity component that is capable of being poisoned. Another difficulty is presented when a hydrocracking process is used to maximize naphtha yields from a feedstock containing light catalytic cycle oil which has a very high aromatics content. The saturation properties of the catalyst must be carefully gauged to saturate only one aromatic ring of a polynuclear aromatic compound such as naphthalene in order to preserve desirable high octane value aromatic-containing hydrocarbons for the naphtha fraction. If the saturation activity is too high, all of the aromatic rings will be saturated and subsequently cracked to lower octane value paraffins.\nOn the other hand, distillate fuels such as diesel fuel or aviation fuel have specifications that stipulate a low aromatic content. This is due to the undesirable smoke production caused by the combustion of aromatics in diesel engines and jet engines.\nPrior art processes designed to convert high nitrogen content feedstocks are usually two stage processes wherein the first stage is designed to convert organic nitrogen compounds to ammonia prior to contracting with a hydrocracking catalyst which contained a high amount of cracking component; i.e., a molecular sieve material.\nAccordingly, it is an object of the present invention to provide a process for hydrocracking a hydrocracking feedstock.\nMore particularly, it is an object of the present invention to provide a process for upgrading crude shale oil to distillate fuels.\nOther objects and advantages of the present invention will be apparent to those skilled in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to an auto focus system having an AF frame auto-tracking function, and more particularly, to an auto focus system having an AF frame auto-tracking function that allows an AF area (AF frame), which is the target range of an object brought into focus by the auto focus (AF), to follow a predetermined object moved on an imaging screen.\n2. Description of the Related Art\nBroadcasting or commercial television camera systems having an auto focus function for auto focusing on a desired object have been known. This type of camera system generally uses contrast-type AF. In the contrast-type AF, the camera system detects the level of the contrast of the captured image, and controls the focus of an imaging optical system such that the contrast becomes the maximum (the largest).\nIn general, the target range of an object that is brought into focus by AF is not the entire screen of the captured image, but is limited to a portion of the screen, which is called an AF area. Focus control is performed such that the contrast of the captured image (object image) in the AF area becomes the maximum, thereby focusing on the object in the AF area.\nIn the specification, a frame indicating the outline of the AF area is referred to as an AF frame.\nJP-A-2006-258944 (corresponding to US-A-2006/210260) discloses an auto focus system having an AF frame auto-tracking function that allows the AF frame to automatically follow a desired object moved on the screen of the captured image such that the object is continuously in focus. The auto-tracking of the AF frame is performed by detecting the image of an object to be tracked from the captured image and moving the AF frame to the detected position. As a method of detecting the image of an object to be tracked, JP-A-2006-258944 discloses a method of storing the image of the object to be tracked as a reference pattern and detecting an image matched with the reference pattern from the captured image using a pattern matching process.\nIn addition, the following method has been proposed: when the object to be tracked is a person's face, a known face detecting process is performed to detect a face image as the image of the object to be tracked.\nHowever, during the tracking of the face image detected by the face detecting process, when the face is turned 30 degrees or more from the front side in the vertical or horizontal direction, the object may not be recognized as the face image, and tracking may be unavailable.\nWhen the pattern matching process is performed to track a predetermined object, it takes a long time to perform the process, and it may be difficult to track an object being moved at a high speed."}
{"text": "This invention relates to a conveyor for automatic processing of workpieces. One specific aspect by way of example and not limitation relates to a conveyor for the automatic heating, in an oven relying mainly on forced convection, of many square thin flat chips of polypropylene to a suitable temperature in connection with which the invention will be described.\nAn object of this invention is to provide means for conveying flat workpieces on removable trays on an endless type of conveyor through one or more processes.\nIn summary, the invention embraces a conveyor of the endless type having means for clamping and releasing carrier trays that comprises--a pair of parallel and spaced apart endless conveyors such as belts or chains and preferably (as described e.g. hereafter) chains, a plurality of annular rectangular frames connected to and between said chains at equally spaced-apart distances along each and having substantially flat upwardly facing support surfaces; at least one clamping means mounted on each said frame to clamp and release a carrier tray means thereto, a machine base from which said chains are supported to run longitudinally each in an endless path; and a clamp actuator in the form of a fixed cam mounted on said base at a position to engage each of said clamping means seriatim to cause it to release its tray at one position and to clamp its respective tray at another position, said clamping means and clamp actuators coacting to hold the clamping means in an unclamped position over a certain part of the path of said chains and to cause them to move to a clamped position during the other part of the path. In the preferred embodiment, a plurality of such carrier trays and each said carrier tray is a flat perforated member silicone rubber covered on at least one side."}
{"text": "1. Field of the Invention\nThe present invention generally relates to optical lenses and projection systems, and more particularly to an achromatic, Fresnel lens which may be used in a low-profile overhead projector.\n2. Description of the Prior Art\nTransmissive overhead projectors (OHP's) are known in the art, and are generally comprised of a base having a transparent stage area, a light source inside the base, a projection head mounted above the stage, and a condensing lens system located near the stage for collecting and directing the light towards the projector head. The condensing lens system often takes the form of a Fresnel lens or a two-element Fresnel lens combination, as depicted in U.S. Pat. No. 4,436,393.\nSeveral attempts have been made to decrease the base height of OHP's, for portability and reduced obstruction of the projected image during use. To achieve such a low-profile base, some machines collapse during storage or transportation, then expand open for actual use. Machines of this type are described, for example, in U.S. Pat. Nos. 4,969,733 and German Patent Application No. 3,519,506. These collapsible OHP's have several disadvantages, however, such as requiring additional parts which increase the complexity of the machine and add to its cost; the moving parts may also adversely affect alignment of the optical elements over time. Finally, there is no base height advantage of these machines over conventional OHP's when in operation.\nOther attempts have been made to construct a low base height OHP without the need to collapse the base and optical elements. For example, U.S. Pat. No. 4,741,613 employs a three-element refracting Fresnel lens system to allow the point light source to be very close to the stage. There are several aspects of this invention, however, which could be improved. First, the requirement of three Fresnel lenses necessarily adds to the cost. Secondly, this system does not adequately compensate for the diminished illumination near the periphery of the stage area. Finally, the dispersion (chromaticity) of a conventional Fresnel lens may lead to coloration of the projected image, particularly around the periphery thereof. Most OHP's locate the light source within the base to minimize this effect, but it is noticeable any time the projection head is moved from its optimum position. For example, if the projection head must be lowered to magnify an image or focus it on the projection screen, less red light at the periphery of the stage is captured by the projection lens, leading to a bluish tint at the border of the projected image. Similarly, moving the projection head up (away from the stage) can cause the border to have a reddish tint. This effect (which is even more pronounced in the three-element system of the '613 patent) can be overcome by repositioning the lamp within the base (referred to as \"color tuning\"), but this involves further user manipulation of the OHP and still requires a subjective judgment regarding the chromaticity of the projected image. Color tunability also adds to the cost of the OHP since it requires a mechanism for repositioning the lamp.\nThe OHP's disclosed in U.S. Pat. Nos. 4,118,761 and 4,952,053 use folded optical paths to provide a more compact base. In the '761 patent, the light source is also \"off-axis,\" meaning that the apparent location of the light source does not coincide with the center of the stage, i.e., the apparent location is displaced from the normal to the stage center. That device requires a complex optical system including a parabolic reflector to provide collimated light, two planar grooved reflectors, and the condensing lens assembly; it may exhibit a slightly darkened edge, furthest from the light source. In the '053 patent, the darkened edges are compensated for by providing two light sources. The twisted and curved mirrored, grooved prisms are also more difficult to fabricate than flat condensers. Finally, folded optical paths may create problems with full-size stage formats, such as European (A4) styles, since extremely tight folds such as are necessary to achieve a low-profile may cause the light source itself to interfere with the folded light path.\nThe OHP of the present invention overcomes several of the above limitations by using a novel catadioptric Fresnel lens. While the use of any catadioptric lens in an OHP is in itself novel, catadioptric lenses are known in other art areas. The term \"catadioptric\" refers to a lens which uses both reflection and refraction to redirect or bend light waves. See, e.g., U.S. Pat. Nos. 2,003,804, 4,859,043 and 5,097,395. As depicted in those patents, catadioptric lenses are useful in collimating light, similar to Fresnel lenses. In U.S. Pat. No. 5,150,966, a catadioptric lens is used to adjust light intensity distribution by discarding light. See also U.S. Pat. No. 4,755,921, which describes a single element catadioptric lens having increased efficiency at high entrance angles, thereby providing a \"fast\" lens (low f-number) that is suited for compact optical devices.\nChromatic aberration is also a problem with catadioptric Fresnel lenses. The dispersion is caused by the optical properties of the lens material, i.e., its index of refraction varies with the wavelength of the light passing therethrough. There are several techniques, however, for minimizing this effect. For example, attempts have also been made to minimize chromatic aberration in single element lenses, such as by providing discontinuous axisymmetric surfaces, or by using aspheric surfaces. Diffraction gratings may be placed on the lens, including on a Fresnel lens, as depicted in U.S. Pat. No. 5,161,057. See also U.S. Pat. No. 4,337,759, which describes a curved base, catadioptric Fresnel lens with a controlled amount of chromatic dispersion for use as a solar concentrator. Many of the foregoing techniques minimize dispersion only for selected wavelengths, as opposed to the full (visible) spectrum. Most of the techniques (including the '759 patent) also require materials having specific indices of refraction, to achieve a specific focal length with color correction.\nIn light of the foregoing, it would be desirable and innovative to design a Fresnel lens which is usable in a low-profile OHP. The lens would advantageously be achromatic and constructed of any number of light-transmissive materials having a large range of refractive indices, and efficiently focus the light to the projection lens of the OHP. The condensing system preferably would avoid the use of micro-precision diffraction gratings, and would generally minimize the number of optical elements (such as grooved mirrors), including the number of elements in the condensing lens, but still be compatible with a folded optical path. Provision of such an achromatic lens would greatly reduce or minimize the need for color tuning."}
{"text": "1. Field of the Invention\nThe present invention relates to a hoisting apparatus for a load such as luminaires used at high elevations, and particularly a hoisting apparatus with a horizontal stabilizing means for a load holder, to which the load is detachably attached.\n2. Disclosure of the Prior Art\nIn high-ceilinged structures such as concert hall, gymnasium, and convention hall, a hoisting apparatus has been utilized to readily carry out maintenance work of luminaires operated at high elevations. For example, as shown in FIG. 10, European Patent Publication No. 1193442 A2 discloses this kind of hoisting apparatus. That is, this hoisting apparatus 1S is mainly composed of a base 2S secured to a ceiling, a load holder 3S, to which a luminaire is detachably attached, a pair of cables 4S extending between the base and the load holder, a drive unit (not shown) mounted to the base, to which one ends of the cables are connected, so that the drive unit is operative to take in or let out the cables from the base.\nFrom the viewpoint of facilitating the maintenance work of the luminaire under safe working condition, this hoisting apparatus also has a cable-length adjusting means in the load holder 3S, by which the cable length can be readily adjusted such that a descending position of the load holder matches a position adequate for the maintenance work of the luminaire to avoid dangerous operations at high elevations,\nAs shown in FIGS. 11A to 11C, the cable-length adjusting means comprises a winding shaft 60S of a round-bar shape, which is rotatably supported in a holder case 30S of the load holder 3S. Each of the cables 4S is fixed connected at its one end to a winding drum (not shown) of the drive unit and at the opposite end to the winding shaft 60S. The winding shaft 60S also has an elongate hole 64S, to which a dedicated tool 48S can be inserted through a slit 34S formed in the upper surface of the holder case 30S to inhibit the rotation of the winding shaft. In addition, the winding shaft 60S has engagement grooves 65S at its opposite ends, to which a tip of the dedicated tool 48S can be engaged. The grooves 65S of the winding shaft 60S are exposed to be accessible from outside of the holder case 30S, as shown in FIG. 10.\nBy use of the cable-length adjusting means described above, the cable length can be adjusted as follows. That is, the tip of the dedicated tool 48S is engaged to one of the grooves 65S of the winding shaft 60S, and then the winding shaft is rotated by the dedicated tool to wind desired amounts of the cables thereon. In this hoisting apparatus, since both of the cables 4S are wound around the winding shaft 60S in the same winding direction, as shown in FIG. 11A, the winding operation of the cables can be achieved by rotating the winding shaft 60S. In addition, since a proper tension is applied to the cables 4S under the suspended condition of the load holder 3S, it is possible to readily wind the cables around the winding shaft 60S without looseness. After the desired amounts of the cables 4S are wound on the winding shaft 60S, the dedicated tool 48S is removed from the groove 65S, and inserted into the elongate hole 64S of the winding shaft through the slit 34S of the holder case 30S, so that the rotation of the winding shaft is inhibited to prevent unwinding of the cables 4S from the winding shaft.\nBy the way, from another viewpoint of preventing an inclination of the load holder with the luminaire, an improvement of the conventional hoisting apparatus is being awaited. For example, due to variations in cable length, and variations in position of the cable ends respectively connected to the winding shaft and the winding drum of the drive unit, it is difficult to provide the same effective length between the pair of cables. This means that the inclination of the load holder may easily occurs.\nAccording to the hoisting apparatus described above, it is possible to adjust the cable length by rotating the winding shaft to obtain the desired effective length of the cables between the base and the load holder. However, since both of the cables are wound around the winding shaft in the same winding direction, this cable-length adjusting means do not have a function of correcting the inclination of the load holder in a horizontal position. In other words, when one of the cables is wound around the winding shaft by rotating the winding shaft, the other one is also wound around the winding shaft. Therefore, the inclination of the load holder can not be solved by the rotation of the winding shaft.\nThus, from the above viewpoint of stably maintaining the load holder in a horizontal position, there is still room for improvement of the conventional hoisting apparatus."}
{"text": "Haematite, having the chemical formula Fe2O3, is one of the most abundant minerals in nature. It exists as iron ore, in other minerals such as bauxite, and is also a component in clay minerals. It is the major component in laeritic soils (red soils found in the tropics). Similarly, manganese oxide, having a formula Mn2O3 is also a very common component in several laeritic soils and also exists as a mineral of manganese in the tropics.\nU.S. Pat. Nos. 5,645,518 and 5,830,815 issued to Wagh et al. on Jul. 8, 1997 and Nov. 3, 1998, respectively, disclose processes for utilizing phosphate ceramics to encapsulate waste. U.S. Pat. No. 5,846,894 issued to Singh et al. on Dec. 8, 1998 discloses a method to produce phosphate bonded structural products from high volume benign wastes. None of these patents provides a method for utilizing the waste materials of iron and manganese.\nU.S. Pat. No. 6,153,809 issued to Singh et al. Nov. 28, 2000 and U.S. patent application Ser. No. 09/751,655 filed Dec. 29, 2000, publication no. U.S. 2002/0123422 to Wagh et al. represent additional development of the use of chemically bonded phosphate ceramics to useful materials. Each of the aforementioned patents, that is U.S. Pat. No. 5,645,518 issued to Wagh et al., U.S. Pat. No. 5,846,894 issued to Singh et al., U.S. Pat. No. 5,830,815 issued to Wagh et al., U.S. Pat. No. 6,153,809 issued to Singh et al., U.S. Pat. No., 6,133,498 issued to Singh et al. and the above-identified publication no. US 2002/0123422 (patent application Ser. No. 09/751,655) is incorporated herein in their entireties.\nThe phosphate ceramics disclosed in the various patents and publication hereinbefore mentioned illustrate a continuing effort to use the chemically bonded phosphate ceramics disclosed therein for a variety of purposes including the encapsulation of hazardous or radioactive waste as seen in the aforementioned publication, as well as the production of low cost structural materials. Accordingly, therefore, a need exists in the art for a low cost structural material which combines with synthetic organic resin based structures, for particular usage in the construction industry. Typically, in warm weather climates, low cost housing may be constructed using styrofoam as a base material onto which is sprayed a concrete-like material as a finish coating to seal the styrofoam base material against the elements and to provide a satisfactory looking structure. Heretofore, the phosphate ceramics disclosed in the above-captioned patents and publication were used as a finish coating in warm temperature climates but have not been satisfactory because the bond between styrofoam and the phosphate ceramics herein above disclosed is physical and peelable such that durable coatings have not been able to be provided with the extant material."}
{"text": "The goal of plant breeding is to combine, in a single variety or hybrid, various desirable traits. For field crops, these traits may include resistance to diseases and insects, resistance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination, stand establishment, growth rate, maturity, plant height and ear height, is important. Traditional plant breeding is an important tool in developing new and improved commercial crops."}
{"text": "Recent development in the image sensing field has resulted in image sensors that feature security and machine vision applications, including gesture sensing, depth analysis, iris detection, eye tracking, night or low light vision, etc. In one aspect, these image sensors use traditional, VIS spectrum image sensing pixels to produce VIS images. In another aspect, these sensors also use additional, IR spectrum pixels to produce IR images. Basically, these sensors serve a dual purpose of producing both VIS and IR images.\nConventionally, such a dual purpose image sensor employs a direct combination design wherein an IR sub-sensor is physically juxtaposed next to a VIS sub-sensor. The VIS sub-sensor only includes VIS imaging pixels, and the IR sub-sensor only includes IR imaging pixels. This design offers simplicity, but the VIS and IR sub-sensors inevitably have different vantage points with regard to an imaging target. As a result, the resulting VIS and IR images not only have different optical spectra (VIS versus IR), but are also not able to be completely superimposed onto each other. This creates difficulties for subsequent image processing.\nCorresponding reference characters indicate corresponding components throughout the several views of the drawings. Skilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present invention. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments of the present invention."}
{"text": "Degradable polyelectrolytes are useful materials in a variety of biological applications ranging from biomedical implant coatings and immunostimulants to vehicles for drug and nucleic acid delivery. A significant fraction of these applications rely on polyelectrolytes that perform a specific function, after which they degrade into non-toxic byproducts, thereby preventing bioaccumulation and toxicity. Many known degradable polyelectrolytes contain hydrolytically unstable functionalities in the polymer backbone, including ester, anhydride, acetal, carbonate, amide, phosphate, and siloxy ether functionalities. While natural polyelectrolytes such as collagen and chitosan have garnered substantial interest in various biological applications, new approaches to degradable synthetic polyelectrolytes will continue to furnish well-defined materials with tunable structures, controlled molecular weights, and variable backbone charge densities and hydrophilicities.\nSome of the most widely studied synthetic degradable polyelectrolytes are based on polyphosphazene and poly(β-amino ester) scaffolds. Anionic and cationic derivatives of poly(phosphazene) are promising vaccine adjuvants and nucleic acid delivery agents, respectively; however, complex monomer syntheses and harsh polymerization conditions limit the types of chemical functionality that may be introduced into these materials. Synthesized by the Michael addition polymerization of diamines with diacrylates, poly(beta-amino esters) comprise a modular platform of polycationic materials exhibiting highly variable hydrophilicities and tunable degradabilities depending on the specific monomers used. Various groups have demonstrated the utility of these materials as components in drug delivery vehicles and in erodible polyelectrolyte multilayer films for therapeutic small molecule and nucleic acid delivery. In spite of the demonstrated potential of poly(beta-amino esters), their widespread utility in biomedical applications is curtailed by synthetic difficulties associated with tuning the charge density along the polymer backbone.\nWhat is needed are readily accessible monomer and polymer platforms that provide new materials, particularly for biological applications."}
{"text": "A biomarker is a measurable characteristic that reflects the severity or presence of or is associated with some disease state and that can be used as an indicator of a particular disease state or some other physiological state of an organism. Biomarkers can be specific cells, molecules, genes, gene products, enzymes, receptors, mutated versions of any of these cellular elements or hormones that can be used to identify and/or measure the presence or progress of disease state, such as a particular cancer or tumor. Further, it is well known that tumors and cancers can express a set of tumor biomarkers that can be used to identify the presence of or measure the progress of or the effects of treatment on the tumor or cancer.\nDespite the abundant use of biomarkers for diagnosing disease and monitoring progression of the same, there remains a need for developing therapeutic approaches that make use of this information to specifically target biomarkers expressed by the disease that are directly associated with the proliferation or existence of the diseased state and subsequent deterioration of a subject's overall health\nListeria monocytogenes (Lm) is an intracellular pathogen that primarily infects antigen presenting cells and has adapted for life in the cytoplasm of these cells. Listeria monocytogenes and a protein it produces named listeriolysin O (LLO) have strong adjuvant properties that unlike the majority of adjuvants used for cellular based immunotherapies, can be administered after providing an antigen specific treatment or can be used to itself provide antigen-specific treatment when fusing an antigen of interest to an adjuvant protein expressed by the Listeria, such as LLO or an ActA protein.\nThe present invention addresses this need by providing a combinatorial, multi-target immunotherapeutic approach wherein individual compositions each comprising a recombinant Listeria-strain expressing a different disease-associated antigen than a counterpart Listeria present in a separate composition, are administered separately to a subject having a disease, or the compositions are administered in combination as single bolus administration. The present invention further addresses this need by providing a predetermined number disease-associated antigens or fragments thereof by using a recombinant Listeria expressing at least one fusion protein comprising the antigen fused to an immunogenic Listeria peptide such as an N-terminal LLO, truncated LLO, an ActA protein fragment, or a PEST peptide. Use of such compositions will allow diseases, including tumors, cancers, or others having sub-populations of diseased cells expressing more than one biomarker to be successfully treated."}
{"text": "Surface representations in computer graphics can be orders of magnitude smaller than polygonal or patch based surface representations which may be a desirable feature for memory constrained devices like game consoles. Constructive Solid Geometry (CSG) operations are a powerful method for defining surfaces of high genus (e.g. a sphere with a hole) from surfaces of low genus, such as plain cylinders and spheres. Defining parametric procedural surfaces of genus 0 (e.g., a sphere) and some surfaces of genus 1 (e.g., a torus) is usually straightforward since there are no holes in the surface domain. However, for surfaces of higher genus, it is much more difficult since domains of some parametric surfaces may have holes whose boundaries are, in some sense, only implicitly defined (e.g., as an intersection of surfaces). These boundaries are difficult to program manually.\nCSG operations may provide methods for defining surfaces of high genus from surfaces of low genus, such as cylinders and spheres. Virtually any manufactured object may be modeled using CSG operations in combination with surfaces of revolution and generalized extrusions, both of which are easily programmed procedurally. Also, the addition of CSG operations to procedural surfaces dramatically increases the scope of objects that can be modeled procedurally.\nAmong other things, CSG operations may require computing and representing curves of intersection between the two or more surfaces being operated upon. These curves, in general, are defined only implicitly as the solution to an equation of the form f(x1, . . . , xn)=0. These equations are not easy to evaluate and typically require a sophisticated, slow, and computationally costly global zero finding solver. A more compact, exact and resolution independent representation of such curves on the other hand may be efficiently evaluated at runtime on a graphics processor. Such representation may be a highly desirable feature for memory constrained devices, like game consoles or for bandwidth constrained applications. Also, such representations may make it possible to represent high genus surfaces in an entirely procedural way without significant computational costs."}
{"text": "(1) Field of the Invention\nThe present invention relates to a television device for providing images with appropriate quality in accordance with a viewing environment.\n(2) Description of the Related Art\nVarious functions of the video display device such as a television device have been proposed for adjusting the video image quality in accordance with the viewing environment such as ambient brightness. For example, Japanese Published Unexamined Patent Application No. 4-342373 discloses the television receiver structured to automatically adjust the image quality to the level adapted to the viewing environment based on the illuminance data of the ambient brightness and the distance data with respect to the viewing distance between the receiver and the viewer. Japanese Published Unexamined Patent Application No. 2000-112021 discloses a display device of projection type structured to allow the light detection unit to detect the light flux from the display area, and automatically adjust the light flux projected to the display area based on the output signal from the light detection unit."}
{"text": "The following publications are believed to represent the current state of the art:\nU.S. Pat. Nos. 5,659,771; 5,907,839; 6,424,983; 7,296,019; 5,956,739 and 4,674,065\nU.S. Published Patent Application Nos. 2006/0247914 and 2007/0106937;"}
{"text": "The invention relates to a tool-driving device that is particularly designed for use in machine tools or in machining units of machining centers, and has at least one machine spindle that is seated to move.\nMachine tools are used especially for material-removal processes, such as boring, milling, turning on a lathe, etc.\nThe tool is inserted into a corresponding tool receptacle that is secured in the work spindle of the relevant machine tool. Various tool receptacles are available.\nDuring the machining process, the work spindles are driven by associated drive apparatuses. Control devices, which can include expanded electronic circuits or execution programs, are provided for controlling the spindle movement, notably its rotation and/or adjustment.\nThe control device establishes the rpm of the spindle within an rpm range. This range is inherently limited. It may be that, particularly for very small tool diameters or for other reasons, rpms outside of the rpm range of the spindle are required.\nIt is the object of the invention to provide a tool-driving device that expands the application range of a machine tool or machining unit, preferably with as little intrusion as possible into the existing machine control.\nThis object is accomplished with a tool-driving device having the features of claim 1.\nThe tool-driving device of the invention has a spindle insert, which can preferably be clamped, fixed against relative rotation, in a machine spindle and can support a tool for machining workpieces. A coupling device serves to secure the spindle insert in the machine spindle. A drive that is supplied by a drive source located outside of the spindle insert, and can be controlled by a control device, is provided for driving the tool. The drive is effected by way of a coupling element that can be connected to the supply lines of the drive. The drive is controlled as a function of the movement of the machine spindle; the tool-driving device is provided with a detection device for detecting this movement.\nFrom the spindle movement, the detection device obtains a signal that characterizes, for example, the rpm, and is used as an input signal for the control device for controlling the drive, and therefore the movement, of the tool. The detection of the rpm requires no access to the machine control, especially if no control signals originating from the machine control are necessary. The control device is separate from the other machine control, and is therefore independent and self-sufficient.\nIf desired, the power supply can be effected by the tapping of the machine control or the drive source of the machine tool. A dedicated drive source can, however, also be provided for the power supply.\nThe tool-driving device permits the increase of the rpm of the machine spindles above and beyond the capabilities of the machine spindle. Unlike in a passive accessory gear, in this instance the additional supply of power in the drive of the tool permits the conversion of an output that exceeds the output of the machine spindle. The maximum torque can be completely retained while the rpm is increased.\nThe spindle insert has a coupling device, e.g., a 7/24 taper shank, which permits a secure, detachable connectionxe2x80x94fixed against relative rotationxe2x80x94with the machine spindle. It also has an essentially cylindrical, one- or multiple-part housing, inside which the drive is disposed.\nIf material-removal operations are to be executed with a rotating tool, the drive is embodied as a rotary drive. A motor, e.g., an electric motor, serves to drive the tool. DC motors, synchronous motors or asynchronous motors can be used for a single- or polyphase alternating current. Hydraulic or pneumatic drives, with which rotational or axial movements of the tool can be attained, can also be used. The motors can be connected to the tool directly, or via a gear in a driving arrangement.\nIn a preferred embodiment of the invention, a receiving apparatus is provided for receiving the tool; the apparatus has a tool spindle, into which the tool is clamped, fixed against relative rotation. The tool spindle preferably has a conical inside shape. The tool spindle is then formed by a rotatably-seated shaft, and projects out of the housing. The shaft is connected to a rotating part of the motor (internal or external rotor) so as to be fixed against relative rotation. The shaft and the tool spindle are preferably embodied to rotate symmetrically relative to an axis of rotation established by the machine spindle. The tool spindle can, however, also support a quick-clamping element, a jaw chuck or the like.\nAt least one slip ring, which is mounted to the outside of the housing and is electrically insulated from it, and can be brought into engagement with an associated sliding contact of the coupling element, is provided for supplying power to the electric motor. When the machine spindle rotates, the sliding contacts slide along the slip rings, thereby assuring the power supply to the drive. Rollers can also be used instead of sliding contacts. The supply can also be effected contactless, e.g., with transformers.\nThe slip rings are preferably disposed on a conical part of the housing whose diameter increases starting from the machine spindle. The slip rings therefore have different diameters. The smallest diameter is larger than that of an arbitrary part of the coupling device. Thus, the spindle insert can be inserted into the machine spindle without altering the position of the contact set. The contact set can then be rigidly secured to the machine tool, in which case it is disposed at a slight incline, corresponding to the incline of the conical housing part. The insertion of the spindle insert produces the contact between the slip rings and the sliding contacts. The contact set can also be seated to be adjusted, and/or can be separate.\nThe safety of the tool-driving device is increased when voltage is only applied to the sliding contacts during the machining process. If the detection device detects rpms that are at least as high as a defined threshold value, preferably 30 rpm, the current supply to the sliding contacts is enabled, for example, by the automatic closure of a switch. The circuit is opened at rpms below the threshold value.\nContactless, magnetic or optical methods are preferred for rpm detection. For example, a metal part connected to rotate with the spindle insert or the machine spindle can serve to induce a short voltage pulse in a stationary coil with each rotation.\nIn an advantageous embodiment, the detection device has a signal generator, particularly a light source, and a signal receiver, particularly a light sensor. The detection device is preferably adjustably mounted to the machine tool, for example to the spindle head that guides the machine spindle. A marking, such as a narrow metal plate, that reflects the light emitted by the light source is secured to the tool coupling or the machine spindle. A signal that is thereby generated, and characterizes the rpm of the machine spindle, e.g., a pulse signal that is proportional thereto, is then transmitted to the control device.\nA circular clamping body having different visual properties from the location where it is to be secured can serve as a marking. The clamping body can have a gap or a recess.\nThe passage of the gap or recess in front of the sensor generates the signal.\nMarkings that effect the generation of a plurality of signals with each rotation can also be provided. In the simplest case, the markings can be equidistantly spaced and provided on, for example, an adhesive strip.\nThe control device utilizes the signals arriving from the detection device to generate a corresponding drive signal for the drive. Hence, the rpm range of the tool can be expanded with the device of the invention. Existing machine tools can therefore be rendered more versatile without its mechanical or electronic components being disturbed.\nThe control device can be integrated into the spindle insert, or accommodated separately. It can also be controlled by programs running on a computer. A console can be provided for the user.\nAt least one supply line for a cooling fluid or compressed air is preferably provided in the tool-driving device for cooling the tool, as is an outward-oriented nozzle, which is preferably pivotable and comprises plastic, for example. At the same time, the nozzle can conduct heat out of the tool-driving device.\nFurther advantageous details about embodiments of the invention ensue from the dependent claims, the drawing and/or the associated description."}
{"text": "(1) Field of the Invention\nThe present invention relates to a transmission system having a plurality of terminals linked by a working channel line and a protection channel line in a redundant structure wherein a connection is switched from the working channel line to the protection channel line while avoiding instantaneous cutoff when the working channel line has failed.\nIt is desired to provide reliable data communication in a transmission system having a plurality of terminals linked by a working channel line and a protection channel line in a redundant structure. In order to attain this objective, it is necessary that the transmission system reliably switch a connection from the working channel line to the protection channel line while avoiding instantaneous cutoff when one of the two channel lines has failed.\n(2) Description of the Related Art\nFIG. 22 shows a conventional transmission system having a plurality of terminals linked by a working channel line and a protection channel line in a redundant structure. In the conventional transmission system, a connection is switched from the working channel line to the protection channel line while avoiding instantaneous cutoff when a switching command to switch one of the channel lines to the other is received.\nAs shown in FIG. 22, the conventional transmission system includes a transmission terminal 101 (which is called a terminal A), a reception terminal 103 (which is called a terminal B), and a transport terminal 102 (which is called a terminal C). The terminal A and the terminal B are linked by a working channel line, and the terminal A, the terminal C and the terminal B are linked by a protection channel line. In the conventional transmission system of FIG. 22, transmission of a digital signal in the existing synchronous digital hierarchy (SDH) frame format is assumed.\nIn the transmission terminal 101 (the terminal A), a frame pulse insertion unit (FP INS) 104 and a distribution unit (DIS) 105 are provided. The FP INS 104 inputs the SDH frame and inserts a frame pulse (FP) in the SDH frame at a given location of the SDH frame. The frame pulse (FP) inserted by the FP INS 104 indicates a reference position in the SDH frame on the related channel line. When a switching command is externally supplied to the reception terminal 103, the FP in the SDH frame is used by the reception terminal 103 to synchronize the SDH frame on the working channel line with the SDH frame on the protection channel line.\nThe DIS 105 supplies the SDH frame (with the FP inserted) from the FP INS 104 to both the terminal B through the working channel line and the terminal C through the protection channel line.\nIn the transport terminal 102 (the terminal C), a data delay unit 106 is provided. The data delay unit 106 provides a delay for the SDH frame on the protection channel line from the terminal A.\nIn the reception terminal 103 (the terminal B), a fixed delay unit 107, a frame pulse detection unit (FP DET) 108, and a control unit 109 are provided. The fixed delay unit 107 provides a fixed delay for the SDH frame on the working channel line from the terminal A. The fixed delay is provided by the fixed delay unit 107 such that a total delay for the SDH frame on the working channel line due to the transmission between the terminal A and the terminal B and due to the transmission through the fixed delay unit 107 within the terminal B is always greater than a delay for the SDH frame on the protection channel line due to the transmission between the terminals A, C and B and due to the transmission through the data delay unit 106 within the terminal C.\nThe frame pulse detection unit (FP DET) 108 detects the frame pulse (FP) in the SDH frame on the working channel line, and separates the frame pulse (FP) from the SDH frame on the working channel line. The FP DET 108 supplies the frame pulse to the control unit 109 and outputs the reconstructed SDH frame (with no frame pulse) which is the same as the SDH frame originally sent on the working channel line from the terminal A.\nFurther, in the reception terminal 103, a data delay unit 110, a frame pulse detection unit (FP DET) 111, and a selector 112 are provided. The data delay unit 110 provides a variable delay for the SDH frame on the protection channel line from the terminal C. The delay provided for the SDH frame on the protection channel line by the data delay unit 110 is controlled by the control unit 109 such that a position of the FP in the SDH frame on the protection channel line matches with a position of the FP in the SDH frame on the working channel line.\nThe frame pulse detection unit (FP DET) 111 detects the frame pulse (FP) in the SDH frame on the protection channel line, and separates the frame pulse (FP) from the SDH frame on the protection channel line. The FP DET 111 supplies the frame pulse (FP) to the control unit 109 and outputs the reconstructed SDH frame (with no frame pulse) which is the same as the SDH frame originally sent on the protection channel line from the terminal A.\nThe control unit 109 receives the frame pulse (FP) from the FP DET 108 and the frame pulse (FP) from the FP DET 111, and controls the variable delay of the data delay unit 110 based on the FP from the FP DET 108 such that a position of the FP in the SDH frame on the protection channel line matches with a position of the FP in the SDH frame on the working channel line. Further, the control unit 109 controls the selector 112 in response to an externally supplied switching command, so that the selector 112 outputs a selected one of the SDH frame on the working channel line from the FP DET 108 and the SDH frame on the protection channel line from the FP DET 111.\nThe selector 112 outputs the selected one of the SDH frame on the working channel line from the FP DET 108 and the SDH frame on the protection channel line from the FP DET 111, under the control of the control unit 109, as the output data from the terminal B.\nFIG. 23 shows an operation of the conventional transmission system of FIG. 22.\nIn the above-described conventional system, when a digital signal in the SDH frame format is input to the transmission terminal 101, the frame pulse insertion unit 104 inserts a frame pulse (FP) in the SDH frame. The SDH frame with the FP inserted is produced at the output of the FP INS 104.\nThe distribution unit 105 outputs the SDH frame (with the FP inserted) to both the reception terminal B through the working channel line and the transport terminal C through the protection channel line (\"S1\" in FIG. 23).\nA first delay due to the transmission of the digital signal on the working channel line is provided for the SDH frame having the FP on the working channel line. The delayed SDH frame is input to the fixed delay unit 107 in the terminal B (\"S2\" in FIG. 23).\nA second delay due to the transmission of the digital signal on the protection channel line and through the data delay unit 106 of the terminal C is provided for the SDH frame having the FP on the protection channel line. The delayed SDH frame is input to the data delay unit 110 in the terminal B (\"S3\" in FIG. 23).\nA third delay is further provided for the SDH frame having the FP on the working channel line by the fixed delay unit 107 in the terminal B. The delayed SDH frame is output from the fixed delay unit 107 (\"S4\" in FIG. 23).\nThe third delay (or the fixed delay) is provided by the fixed delay unit 107 such that a total delay (the first delay plus the third delay) for the SDH frame on the working channel line input to the FP DET 108 is always greater than the second delay for the SDH frame on the protection channel line input to the FP DET 111.\nThe control unit 109 controls the variable delay of the data delay unit 110 based on the FP from the FP DET 108 such that a position of the FP in the SDH frame on the protection channel line matches with a position of the FP in the SDH frame on the working channel line. Further, the control unit 109 controls the selector 112 in response to an externally supplied switching command, so that the selector 112 outputs a selected one of the SDH frame on the working channel line from the FP DET 108 and the SDH frame on the protection channel line from the FP DET 111. The selector 112 normally outputs the SDH frame on the working channel line from the FP DET 108. When the switching command to switch the working channel line to the protection channel line is externally supplied to the control unit 109, the selector 112 outputs the SDH frame on the protection channel line from the FP DET 111.\nAccordingly, when both the working channel line and the protection channel line normally function and a switching command to switch one of the channel lines to the other is externally supplied, the conventional transmission system of FIG. 22 can switch a connection from the working channel line to the protection channel line without causing instantaneous cutoff in response to the command. However, when the working channel line has failed, the switching command cannot be supplied to the control unit 109 in the reception terminal 103. Therefore, when the working channel line has failed, it is impossible for the conventional transmission system of FIG. 22 to switch a connection from the working channel line to the protection channel line while avoiding instantaneous cutoff. It is difficult for the conventional transmission system of FIG. 22 to provide reliability for data communication if the working channel line has failed."}
{"text": "1. Field of the Invention\nThe present invention relates to a push button switch for use as a key switch for a data input apparatus for a personal computer, a word processor or the like, and, more particularly, to a push button switch the thickness of which can be easily reduced.\n2. Related Art Statement\nA conventional push button switch of the type described above will be described with reference to FIG. 13.\nReferring to the drawing, reference numeral 1 represents a reinforcing plate made of metal or the like. A membrane switch 2, comprising an upper sheet 2b having a movable contact 2a and a lower sheet 2d having a fixed contact 2c, is placed on the reinforming plate 1. A case 3 is placed on the membrane switch 2, the case 3 having an annular first projecting portion 3a and a second projecting portion 3b. An operation member 4a of a key top 4 is movably positioned along an inner surface 3c of the first projecting portion 3a. The operation member 4a has a fastening claw at the lower end portion thereof so that the fastening claw is fitted within a recessed portion (omitted from illustration) formed in the first projecting portion 3a. As a result, the fastening claw can be moved within the recessed portion and the operation member 4a can be vertically moved along the inner surface 3c. Furthermore, a coil spring 5 is interposed between a flat surface 3d of the case 3 and the lower surface of the key top 4 in such a manner that the coil spring 5 is positioned around the first projecting portion 3a. In addition, another coil spring 6 is positioned in the operation member 4a by press-fitting for the purpose of pressing the movable contact 2a of the membrane switch 2. Although omitted from the illustration in FIG. 13, a recessed portion is formed in the second projecting portion 3b in a direction perpendicular to the direction of the drawing sheet for FIG. 13. As a result, the upward separation of the key top 4 is prevented by fastening the fastening claw provided for the key top 4 to the recessed portion.\nThus, when the key top 4 is depressed against the urging force of the coil spring 6, the outer portion of the operation member 4a is downwards moved along the inner surface 3c. As a result, the lower end portion of the coil spring 5 press-fitted in the operation member 4a presses the upper sheet 2b of the membrane switch 2, causing the movable contact 2a to be brought into contact with the fixed contact 2b. Therefore, the switch is switched on. When the pressure applied to the key top 4 is then released, the original state can be restored by the elastic restoring force of the coil spring 6.\nRecently, there has been a desire for a compact keyboard having a reduced thickness, causing a necessity for reducing the thickness of the push button switch to arise.\nWhen a push button switch having a reduced thickness is constituted, the push button switch must have a certain depressing stroke (3 to 4 mm).\nIt is assumed that the thickness of the key top 4 is b, the distance of the movement of the key top 4 is S, the length of a portion (omitted from illustration) for fastening the case 3 and the key top 4 is a, the thickness of the reinforcing plate 1 and the membrane switch 2 is c and the overall height is expressed by H.\nThen, the length a of the fastening portion can be expressed by a=H-(2S+b+c). In this case, the thickness b of the key top 4, the movement distance S of the key top 4 and the thickness c of the membrane switch 2 and the reinforcing plate 1 become substantially constant depending upon the molding condition and the parts composition. Therefore, there has conventionally been a necessity for the length a of the fastening portion to be shortened at the time of realizing the above-described thickness reduction.\nHowever, if the length a of the fastening portion is shortened, the lower end portion of the operation member 4a of the key top 4 is caught by the inner surface of the first projecting wall 3a of the case 1 when the end portion A of the key top 4 is depressed. As a result, the conventional push button switch cannot be depressed smoothly."}
{"text": "The present disclosure relates generally to medical devices for implantation within a human or animal body for repair of damaged vessels, ducts, or other physiological pathways, and particularly, to prostheses with side branch lumens.\nEndovascular methods have been proposed for treatment of diseases of the aorta such as aortic dissection and aortic aneurysm. Using prostheses, such as stent grafts, to treat aneurysms is common in the medical field. Stent grafts are deployed by accessing a vasculature with a small incision in the skin and guiding a delivery system to the target area. This endoluminal delivery is less invasive and generally preferred over more intrusive forms of surgery. Multiple stent grafts may be implanted using endoluminal delivery to provide a system of interconnected stent grafts. Interconnected stent grafts can be made of fenestrated stent grafts and smaller side branch grafts, including bifurcated components.\nSuch methods have been proposed particularly when the diseased portion of the aorta is adjacent the aorta bifurcation. But when the diseased portion of the aorta is located higher up in the aorta, for example, in the region of the descending aorta adjacent the thoracic arch or in the ascending aorta, endovascular techniques for treating these diseases are somewhat more difficult because of the arched or curved nature of the thoracic arch, the presence of major arteries in the region, and the proximity to the heart.\nFor instance, for treatment of thoracic aortic aneurysms and/or dissections in particular, it is necessary to introduce the stent graft high up in the aorta and in a region of the aorta which is curved and where there can be strong blood flow. Furthermore, in the thoracic aorta there are major branch vessels extending therefrom, such as the brachiocephalic, carotid and/or subclavian arteries. During and/or after treatment of an aneurysm or dissection in the region of the thoracic arch, it is desirable for blood supply to continue to flow to these branch arteries. For this purpose, fenestrations or side branches are provided in a stent graft that is placed in that region, through which side arms or branch extensions may be deployed and extend into the brachiocephalic, carotid and/or subclavian arteries, for example.\nCustom made devices, including scalloped and fenestrated devices, have been used in situations where the arch vessels are compromised and entire coverage of the aortic arch is not required. However, deployment of these devices may be difficult."}
{"text": "1. Field of the Invention\nThe present invention relates to battery frames used when stacking flat cells to form a battery, and a battery in which the frames are used.\n2. Description of the Related Art\nA plurality of single cells is combined in series or parallel to form a battery or a module battery with high power and capacity. For example, it is known that a plurality of flat single cells, each having an electric-power generating element covered with laminate films, is connected to each other in series or parallel to form a battery with high power and capacity (see Japanese Patent Laid-Open Publication No. 2001-256939). A flat single cell in which an electric-power generating element is covered with laminate films is called a laminate cell.\nWhere a battery is formed using laminate cells of this kind, a plurality of laminate cells is located on the same plane, a plurality of sets of these laminate cells on the same plane is further stacked, and then these laminate cells are connected to each other in series or parallel. By stacking the laminate cells like this, a battery with higher power and capacity can be formed.\nWhen assembling this stack-type battery, locating and stacking laminate cells one by one results in poor workability. Therefore, laminate cells are located on a plate-like frame and a plurality of frames where the cells are located is stacked, thus improving workability in assembling a stack-type battery.\nFurther, in order to electrically connect the cells located on neighboring frames in the stacking direction, washers made of a conductive material (conductive washers) may be incorporated in the frames at positions where electrode tabs of the laminate cells lie. To electrically insulate cells, washers made of insulating material (insulating washers) may be incorporated in the frames. When electrically connecting or insulating the laminate cells located on neighboring frames in the stacking direction, the frames, where the laminate cells are mounted, are stacked and the frames or washers are pressed from the top and bottom of the stack. This ensures that the electrode tabs and the washers come into contact with to each other, enabling the neighboring cells to be electrically connected or insulated from each other."}
{"text": "This invention relates to an improvement to tags, the sort of tags that are, for example, attached to hospital fluid bags, such as blood bags, saline solution bags and the like. To be able to identify any one bag and match it to a particular patient or operation, a tag has to be attached since the bag itself may not be tampered with or written upon. The tag is pre-printed with all the information required to identify the patient, operation and contents of the bag it is meant to be attached to.\nSuch tags are well known in the medical field and have been used for quite some time. However, there are drawbacks to the presently available tags. They are supplied either part printed for attaching to the fluid bags they are ordered for, supplied largely blank for writing on, or a self adhesive label is printed and adhered to a luggage style tag. All these forms of identification are then secured to the fluid bag by different methods such as cable tie, nylon attachment or string."}
{"text": "This invention relates to the transmission of pictorial image data. More particularly, it is concerned with the progressive transmission and reconstruction of coded images in which an approximate image is reconstructed based upon partial information and details are added as additional information becomes available.\nThe progressive transmission and reconstruction of coded images allows an approximate image based upon partially received information to be constructed to which additional details are added as additional information becomes available. This procedure has various applications in the field of image communications, such as for interactive picture retrieving, variable-rate video conferencing, and the quick display of freeze-frame image transmission. One relatively simple scheme proposed by Knowlton U.S. Pat. No. 4,222,076 issued Sept. 9, 1980, deals with spatial domain data for the progressive transmission of gray-scale pictures. This approach has the advantages of simplicity in implementation and no coding distortion in the final reconstructed image. However, due to the nature of successive picture subdivision introduced by this method, the number of accumulated bits of information increases exponentially with each interation.\nOther schemes that deal with transform domain data have been described in articles by Takikawa \"Fast Progressive Reconstruction of A Transformed Image,\" IEEE Trans. Inform. Theory, vol. IT-30, pp. 111-117, January 1984 and Ngan \"Image Display Techniques Using the Cosine Transform,\" IEEE Trans. Acoust., Speech, Signal Processing, vol. ASSP-32, No. 1, pp. 173-177, February 1984. Transform image coding is well known for its compression efficiency. Its nature renders it also suitable for efficient progressive transmission and reconstruction since low frequency transform coefficients contain most of the energy of image signals. Thus, a small subset of the transform coefficients is good enough for reconstructing a rough version of the whole image, while the remainder of the transform coefficients allow the receiver to add details to the initially reconstructed picture as they are received. In one such scheme the transform coefficients of each block of image data are considered as arranged in a square lattice and are sent and received in a zig-zag pattern in order from the large through the small variance values. This scheme is described in an article of Tescher and Cox \"An Adaptive Transform Coding Algorithm,\" ICC Conference Records, pp. 47.20-24, 1976. Although the zig-zag technique provides better compression efficiency than other proposed transform domain schemes, it is desirable to further improve the efficiency with which image data can be transmitted, particularly during the first few iterations."}
{"text": "1. Field of the Invention\nThe present invention is directed to an optical measurement system for determination of an object's profile or thickness, and more particularly to such an optical measurement system using two optical heads directing individual light beams to different points on the object's surface to measure distances of these points from a reference plane by triangulation for analyzing the surface profile or the thickness of the object based upon the measured distances of the two points on the object's surface.\n2. Description of the Prior Art\nIn order to obtain a depth or height of a step on the surface of an object or thickness of an object by optical triangulation measurement, it has been proposed to use a pair of optical heads disposed to direct individual light beams to different points on the object's surface for measuring the positions of these points. The distances of these points are processed by triangulation and are analyzed to determine the object profile. For example, when the two heads are disposed to measure the positions of the points spaced along the object's surface for measuring individual perpendicular distances to the surface from a reference plane, the difference of the measured distances gives the height or depth of a step existing between these two points. On the other hand, when the optical heads are disposed on the opposite of the object to measure like perpendicular distances of the positions of two points on the opposite surfaces of the object from a reference plane selected to be within the thickness of the object, the addition of the measured distances gives a thickness of the object at these points.\nIn such optical measurement systems, the optical head is normally designed to have a photo-sensor which receives the light beam reflected on a point on the object's surface and provides an output which varies in proportion to the perpendicular distance of the point from a reference plane selected to be generally parallel to the object's surface. The output from the head is processed in an associated signal processing circuit so as to determine a true distance of the point from the reference plane. In this connection, when the two heads are connected to the individual signal processing circuits, there is a potential problem that the distances measured in these separate processing circuits may include individual deviations or discrepancies due to inherent variations in the circuits, for example, deviations in the temperature characteristics of certain elements consisting the circuits. Since these discrepancies are inherent to the individual circuits, they are difficult to be compensated for in obtaining the step in the object's Surface and the thickness of the object. Thus, no reliable analysis is not expected in this system having two optical heads connected respectively to the individual processing circuits."}
{"text": "1. Field of the Invention\nThe present invention relates to an insulated gate field effect semiconductor device (FET) for use in switching devices, integrated circuits (ICs), and display devices such as liquid crystal displays and the like.\n2. Description of the Prior Art\nInsulated gate FETs heretofore fabricated unexceptionally comprise a semiconductor region in which a source, a channel, and a drain are established. In such insulated gate FETs, the drain is in contact with the channel and the source is in contact with the channel. However, in those types of insulated gate FETs, there have occurred problems such as the reverse current leakage from the drain to the source, and the poor drain voltage resistance.\nMore specifically, as is illustrated in FIG. 2, an insulated gate FET of the type above often suffers problems ascribed to the reverse current leakage. That is, the current which flows reversely from the drain to the source yields a curve indicated with (B) in the FIGURE; typically, although the gate voltage VG-drain current ID relation should result in a curve indicated with (A).\nThis phenomenon is ascribed to the occurrence of a punch-through current. That is, even in a gate voltage at which normally no channel forms, i.e., at a condition well below the threshold voltage Vth, an abrupt increase of the drain current occurs if the voltage applied between the source and the drain surpasses a certain value. The generation of this punch-through current is explained by the influence of the reverse bias at the drain junction which also affects the source junction. Since this punch-through current flows between the source and the drain along a path relatively deep with respect to the channel surface, it is possible to cut off the punch-through current by increasing impurity concentration along this path and thereby setting a high resistance between the source and the drain.\nThe low drain voltage resistance also impairs the output characteristics, as is illustrated in FIG. 3 by the curve (B) which shows the drain current ID against the drain voltage VD. At a voltage below the threshold, typically, the ID-VD curve should have a sharp rise as is shown in FIG. 3, curve (A). The low drain voltage resistance is also ascribed to the punch-through current as explained hereinbefore. If an insulated gate FET having a VD-ID curve (B) in FIG. 3 were to be fabricated, a drain current will flow continuously to result in a throw leakage state even though a voltage well below the threshold voltage were to be applied to the gate electrode. This would result in a switching device having poor reliability and insufficient performance.\nAs a means to overcome the problem of punch-through current attributed to the low drain voltage resistance, i.e., the poor insulation between the source and the drain, there is proposed, as is shown in FIG. 4, to provide a semiconductor layer having added therein hydrogen as an offset gate 49. Referring to FIG. 4, there is provided an insulated gate FET comprising a quartz substrate 41, a thin film of polycrystalline silicon 42, a silicon oxide film 43, a polycrystalline silicon electrode 44, a source 45, a drain 46, aluminum electrodes 47, and an offset gate 49. The offset gate prevents the electric field from concentrating in this portion. There is proposed another measure which comprises establishing, to the same area as that of the offset gate, a drain having lightly doped with an impurity which imparts one conductivity type thereto. This process, which is known as a light-dope drain (LDD) process, also relaxes concentration of the electric field in the boundary between the channel and the gate or between the channel and the source. In this process, however, the impurity having doped for imparting one conductivity type to a part of the semiconductor layer diffuses from the drain and the source that there still remains a problem to be solved.\nThis is because the impurity for imparting one conductivity type to the semiconductor layer is the one that is easily diffused by heat. For an insulated gate field effect transistor in which the channel length is not longer than sub-microns, this becomes a major problem. That is, there is a problem that a current flow is continuously formed between the source and the drain in the channel forming region due to diffusion of impurities from the source and the drain to the channel forming region."}
{"text": "Electrodeposited zinc/iron alloys of a semi-bright to lustrous appearance are desirable to provide a decorative plating appearance while simultaneously imparting excellent corrosion protection. Generally speaking, zinc/iron alloys can be deposited on a conductive substrate by means of a zinc/iron alloy electroplating bath. Such electroplating baths and the processes using them are employed to provide alloy deposits on a variety of substrates and are often used in conjunction with ferrous substrates, such as iron or steel.\nThe zinc/iron alloy plating baths and process of the present invention involve the use of a brightening additive which can be used in a wide variety of types of plating baths over broad pH and current density ranges to provide a semi-bright to bright zinc/iron alloy deposit. The plating baths of the present invention are commercially useful and are characterized, in part, by their flexibility and versatility in use to obtain excellent zinc/iron alloy plating results.\nA further understanding of the present invention will be obtained from the following description and examples thereof. Unless otherwise indicated, in the following description and examples, all parts are percents by weight and all temperatures are in degrees Farenheit."}
{"text": "Conventionally used photo or picture frame, whether elegant or general, is totally unable to give photo or picture viewers any benefit of illumination to help them see it clearly at night. Such inconvenience exists particularly when a photo or picture is decoratively mounted high in the dimness of lights. To solve the problem, an illuminant photo frame is greatly needed."}
{"text": "The present invention relates to a brewing head of an espresso maker. The brewing head includes a pressure sealed brewing chamber, and a movable lifting piston forming a floor of the brewing chamber. The movable lifting piston has apertures therethrough. A fill opening is in communication with the brewing chamber for supplying the brewing chamber with fresh ground espresso. A water inlet is provided for feeding hot water through the apertures and into the brewing chamber. The movable lifting piston lifts a cake of used ground espresso located thereon up to an upper edge of the brewing chamber after brewing the espresso. A clearing element laterally pushes the cake off from the lifting piston after the cake is moved to the upper edge of the brewing chamber by the lifting piston, so as to clear the brewing chamber of used ground espresso.\nSuch an espresso maker is described by European Patent No. 0 443 054. In this known espresso maker, the water supply and the clearing blade are arranged on a carriage which is linearly displaceable above the brewing chamber. The brewing process of this device occurs at a pressure of approximately 12 bar.\nA cake of leached espresso grounds is lifted by the lifting piston, and laterally transported off the lifting piston, by a semicircular-shaped clearing element disposed on the aforementioned carriage.\nThe inlet for the hot water is configured as a flexible hose, which represents a source of danger because of the very high pressure involved.\nThis known espresso maker has several drawbacks. For example, during the brewing process, the brewing chamber and the carriage must be sealed from each other. Seals are used which are located in a region of a sliding plane between the brewing chamber and carriage, and are therefore subjected to heavy wear. Further, the linear, reciprocating movement of the carriage also requires a very complicated mechanical control system and a flexible water supply (i.e., a flexible hose). Moreover, the water supply is also under a pressure of 12 bar, and therefore represents a source of danger."}
{"text": "The present invention relates generally to casters, and more particularly to, casters having springs to absorb shock and vibrations.\nCasters have been employed to protect loads being transported by motorized or non-powered material handling equipment from damage due to rough and uneven surfaces, and obstacles that strike the caster wheels. Shock from obstacles or uneven surfaces may be greatly reduced by using spring loaded casters. The shock dampening effect of spring loaded casters reduces wear on the wheel, caster and conveyance as well. Noise may also be greatly reduced when spring loaded casters are used.\nTypically, a spring loaded caster may be a swivel or rigid caster where the wheel and axle assembly is supported with two or more springs to absorb shock and vibrations caused by uneven surfaces, foreign objects and other irregularities of the wheel-engaging surface. Once the spring loaded caster encounters any interference, the springs are compressed so as to absorb the shock, which protects the load from possible damage or sliding off the conveyance.\nIn one embodiment of the present invention, a pivotable spring loaded caster includes a support plate adapted to be affixed to the underside of an article to be transported. A swivel bearing assembly is attached at a top side to the underside of the support plate. Two caster legs depend downwardly from the rotatable portion of the swivel bearing assembly. The caster legs are spaced apart and preferably substantially parallel with each other. Each of the caster legs includes an elongated slot for receiving an axle. A pair of spring brackets is provided that are of a generally L-shaped configuration and are disposed on an outer surface of a corresponding caster leg. The spring bracket may include a wear plate.\nA spring sleeve is disposed on the underside of one leg of each of the spring brackets. Additionally, two or more spring block assemblies are provided. Each spring block assembly includes an axle-receiving segment that includes a passage for receiving an axle therethrough. On each spring block, a spring is held in place by a spring pin and is received in the corresponding spring sleeve. A wheel is provided in the space between the caster leg and an axle that holds the various parts components together is threaded through the spring block assemblies, caster legs, and the wheel.\nIn an alternate embodiment, the rigid spring loaded caster is provided which lacks the swivel bearing assembly of the first embodiment of the present invention. Instead, the two caster legs depend downwardly from the underside of the support plate. Preferably, in this alternate embodiment, first and second wear plates are provided on each of the spring brackets."}
{"text": "The present invention is a layered electrophotographic photoconductor, i.e., a photoconductor having a metal ground plane member on which a charge generation layer and a charge transport layer are coated, in that order. Although these layers are generally separate from each other, they may be combined into a single layer, which provides both charge generation and charge transport functions. Such a photoconductor may optionally include a barrier layer located between the metal ground plane member and the charge generation layer, and/or an adhesion-promoting layer located between the barrier (or ground plane member) and charge generation layer, and/or an overcoat layer on the top surface of the charge transport layer.\nIn electrophotography, a latent image is created on the surface of an insulating, photoconducting material by selectively exposing an area of this surface to light. A difference in electrostatic charge density is created between the areas on the surface exposed and those unexposed to the light. The latent electrostatic image is developed into a visible image by electrostatic toners containing pigment components and thermoplastic components. The toners, which may be liquids or powders, are selectively attracted to the photoconductor surface, either exposed or unexposed to light, depending upon the relative electrostatic charge on the photoconductor surface and the toner. The photoconductor may be either positively or negatively charged, and the toner system similarly may contain negatively- or positively-charged particles.\nA sheet of paper or intermediate transfer medium is given an electrostatic charge opposite that of the toner and then passed close to the photoconductor's surface, pulling the toner from the photoconductor surface onto the paper or the transfer medium still in the pattern of the image developed from the photoconductor surface. A set of fuser rolls melts and fixes the toner on the paper, subsequent to direct transfer or indirect transfer when an intermediate transfer medium is used, producing the printed image.\nThe electrostatic printing process, therefore, comprises an on-going series of steps in which the photoconductor surface is charged and discharged as the printing takes place. It is important to keep the charge voltage on the surface of the photoconductor relatively constant as different pages are printed to make sure that the quality of the images produced is uniform (cycling stability). If the charge/discharge voltage is changed significantly each time the drum is cycled, i.e., if there is fatigue or other significant change in the photoconductor surface, the quality of the pages printed will not be uniform and will not be satisfactory.\nHydrazone derivatives, which have frequently been employed as charge transfer molecules and organic photoconductors for electrophotography, possess interesting photochemical properties which are known to connect closely with the so-called fatigue phenomenon of photoconductors. A good deal of research supports the fact that photoisomerization and photochemical reactions are responsible in large part for the fatigue phenomenon. For example, p-(diethylamino) benzaldehyde diphenyl hydrazone (DEH) undergoes a photochemically-induced unimolecular rearrangement to the indazole derivative, 1-phenyl-3-(4-(diethylamino)-1-phenyl)-1,3-indazole. The following articles give an overview of the mechanism of photo-induced fatigue in electrophotographic conductors: J. Pacansky, et al., Chem. Mater. 4:401(1992); T. Nakazawa, et al, Chem. Lett. 1992, 1125; and E. Matsuda, et al., Chem. Lett. 1992, 1129.\nIn order to use hydrazones as charge transport molecules for electrophotographic applications, photo-induced fatigue has to be reduced to an acceptable level. There are two major paths to minimize photo-induced chemical changes in hydrazone molecules such that photo-induced fatigue of photoconductors can be improved: (1) introducing appropriate substitution on the hydrazone molecules to increase rigidity such that photo-induced cyclization or isomerization can be hindered; and (2) formulating with additives, e.g., a light absorber, in the charge transport layer to filter away the harmful wavelength light (See, for example, U.S. Pat. No. 4,362,798, Anderson, et al.). The former approach will inevitably increase the cost to produce the molecules as compared with the corresponding unsubstituted hydrazones. Therefore, the approach of current choice is the use of additives, such as Acetosol Yellow, to serve as a light filter. Although this approach is effective in reducing room light fatigue of the photoconductor to a certain degree, it also negatively effects the electrical properties of the photoreceptor by increasing discharge voltage and dark decay.\nAzines, which are the product of condensing the remaining NH.sub.2 of a hydrazone with a carbonyl compound, have been disclosed for use in electrophotographic applications, both as transport molecules and as dopants in charge transport layers. Several series of hydrazones and azines are disclosed as charge transport materials in DE3716982, JP62006262 and JP61209456. In addition, some azines have been taught to be used in combination with hydrazones in electrophotographic conductors (see, for example, JP61043752, JP61043753, and JP61043754). It is important to note that these azines are not the fluorenyl-azine derivatives used in the present invention.\nFluorenyl-azines are known in the art. For example, 9-[p-(diethylamino) benzylidenehydrazono)] fluorene has been disclosed in JP57138644 and JP59195659 as a charge transport agent.\nU.S. Pat. No. 4,415,640, Goto, et al, issued Nov. 15, 1983, discloses flourenyl-azines of the type utilized in the present development. The materials are disclosed as charge transport materials, not as adjunct materials used together with another charge transport molecule (see, for example, column 6, lines 52-54; column 7, lines 30-32; and column 8, lines 62-68). The use of these fluorenyl-azines as charge transport materials is taught to minimize photoconductor fatigue.\nIt has now unexpectedly been found that addition to a DEH-containing charge-transport layer of a flourenyl-azine material provides elimination of room light fatigue and cycling fatigue in the resulting photoconductor. For example, a photoconductor containing a DEH-charge transport layer doped with 2-5% azine, exhibits no fatigue after four hours of fluorescent light exposure, while the same photoconductor containing the standard Acetosol Yellow filtering agent exhibits negative fatigue. Increasing the Acetosol Yellow concentration in the charge transport layer, results in negative affects on the sensitivity of the photoconductor and dark decay, while no such effects are observed with the azine material."}
{"text": "1. Field of the Invention\nThis application relates to materials handling. In particular, this application relates to a system, method, and apparatus for tracking items in transit.\n2. Description of the Related Art\nCurrently, when items are in transit between an origin and a destination location, external marking is generally placed on the items to provide information about the item. The information provided about the item by the external marking may be related to the contents or other attributes of the item (e.g., routing, type of handling necessary, etc.). These external markings and other documents are often enclosed in shipping envelopes that are affixed to the item. These shipping envelopes are typically applied by hand, and require a one-time labor effort performed the first time the item is handled during transit.\nRecently, in order to improve the ability to track the location of items, newer tracking technologies have been developed. One of these technologies is radio frequency identification (RFID). RFID technology uses radio waves to obtain information regarding objects involved in the transit process. Electronic tags that carry unique identification and descriptive information are embedded in objects. These tags emit low-power radio frequency signals to RFID readers. RFID readers read RFID tags to obtain the information programmed within the tag's microchip. Readers emit electromagnetic waves from their antennas. Like shipping envelopes, RFID tags are typically affixed to units at the time of shipping, requiring another separate labor effort performed when the item begins its journey.\nThus, items having both shipping envelopes and RFID tags typically require an adhesive for both the shipping envelope and the RFID tag, and they further require two separate labor efforts to affix them both to the item. This duplication of effort and materials results in increased labor and materials costs. Thus, it would be useful to provide a materials handling solution that allows for the affixation of both RFID tags and shipping envelopes without an increase in materials cost or labor."}
{"text": "The present invention generally relates to fuel-fired water heaters and, in representatively illustrated embodiments thereof provides a specially designed high efficiency downfired gas water heater.\nFuel-fired storage type water heaters are commonly used in both commercial and residential applications to provide on-demand hot water to various types of hot water-utilizing plumbing fixtures such as sinks, showers, dishwashers and the like. In one conventional construction thereof, this type of water heater has a tank for holding pressurized water to be heated, a combustion chamber with a fuel burner therein for generating hot combustion products, and a flue extending through the tank interior. During firing of the water heater, hot combustion gases generated by the burner flow through the flue, with heat from the combustion gases being transferred from within the flue to stored tank water through which the flue extends.\nWith increasing demands for both higher energy efficiency and lowered water heater production costs, it has become necessary to design fuel-fired water heaters which are both simpler in structure and capable of transferring a greater percentage of burner-generated combustion heat to the stored tank water. It is to these design goals that the present invention is primarily directed."}
{"text": "1. Field of Invention\nThis invention relates to high voltage power distribution systems and more particularly to apparatus and methods for monitoring for a fault condition in such systems\n2. Description of Related Art\nPower distribution transformers in a high voltage power distribution system are susceptible to internal and external faults. Internal faults may occur due to gradual degradation and/or sudden breakdown of insulating properties of internal structures of the transformer resulting from thermal, electrical and mechanical stresses. The loss of insulating properties of internal structures can lead to a short circuit within the transformer, creating an internal arcing fault. Such fault may occur in a winding of the transformer, for example. A short circuit of this type usually causes burning of the winding, ultimately resulting in an open circuit in the primary winding, whether the fault was initiated in the primary or secondary winding. When such a fault occurs, a fuse, normally installed in series between the high voltage line and the transformer, blows, interrupting current flow to the transformer. Attempts to re-energize the transformer by replacing the fuse can result in further arcing within the transformer and this can rapidly increase the pressure inside the transformer to the point of explosion where burning oil and shards of th transformer casing may be rapidly expelled at the risk of injuring persons and property nearby.\nExternal faults may occur due to insulation failures in a secondary circuit connected to the secondary winding of the transformer. External faults are normally detectable by inspection of the secondary circuit. Re-energizing a transformer having a fault in the secondary circuit do s not normally cause an increase in pressure inside the transformer sufficient to create an explosive threat. Internal faults are thus normally more dangerous than external faults.\nVarious mechanical devices have been devised to detect internal faults, most of which rely on pressure build-up caused by arcing to activate a mechanism. Such devices are described in U.S. Pat. Nos. 5,078,078, 5,623,891, and 6,429,662, for example. Each of these devices employs an efficient but relatively elaborate mechanism to operate a signal device in response to a pressure surge in a transformer caused by an internal fault."}
{"text": "Syndiotactic 1,2-polybutadiene (SPBD) is a thermoplastic resin which can be utilized in making films, fibers and molded articles. For example, U.S. Pat. No. 4,394,473 and U.S. Pat. No. 4,957,970 disclose the use of SPBD in making bags and packaging. It can also be blended into elastomers, such as polydiene rubbers. Because SPBD contains double bonds which are attached in an alternating fashion to its backbone, it can be cocured with the rubbers in such blends. In fact, SPBD/rubber blends provide a unique combination of properties which make them useful in various tire compounds.\nU.S. Pat. No. 4,790,365 discloses that incorporation of SPBD into rubber compositions which are utilized in the supporting carcass or innerliner of tires greatly improves the green strength of those compositions. Electron beam precure (microwave precure) is a technique which has gained wide commercial acceptance as a means of improving the green strength of synthetic elastomers which are used in building tires. However, electron beam precure techniques are costly. The incorporation of SPBD into blends of such synthetic elastomers can often improve green strength to the degree that electron beam precure is not required. The incorporation of SPBD into halogenated butyl rubbers which are utilized as the innerliner compositions for tires also greatly improves the scorch safety of such compositions. U.S. Pat. No. 4,274,462 disclosed that pneumatic tires having improved resistance against heat build-up can be prepared by utilizing SPBD fibers in their tread base rubber.\nAccording to U.S. Pat. No. 4,790,365, the SPBD utilized in making the supporting carcass for tires has a melting point which is within the range of 120.degree. C. to 190.degree. C. and that it is preferred for the SPBD utilized in making the supporting carcass to have a melting point which is within the range of 150.degree. C. to 165.degree. C. The SPBD utilized in making tire innerliners has a melting point which is within the range of 120.degree. C. to 160.degree. C. and preferably has a melting point which is within the range of 125.degree. C. to 150.degree. C. The melting points referred to herein are minimum endotherm values determined from DSC (differential scanning calorimetry) curves.\nTechniques for preparing SPBD by polymerizing 1,3-butadiene monomer are well known in the art. These techniques include solution polymerization, suspension polymerization and emulsion polymerization. The SPBD made utilizing these techniques typically have a melting point within the range of about 195.degree. C. to about 215.degree. C. It is accordingly necessary to reduce the melting point of the SPBD to render it suitable for utilization in some applications.\nA process is disclosed in U.S. Pat. No. 3,778,424 for the preparation of syndiotactic 1,2-polybutadiene which comprises polymerizing 1,3-butadiene in an organic solvent in the presence of a catalyst composition composed of:\n(a) a cobalt compound,\n(b) an organoaluminum compound of the formula AlR.sub.3, in which R is a hydrocarbon radical of 1-6 carbons, and\n(c) carbon disulfide.\nU.S. Pat. No. 3,901,868 reveals a process for producing a butadiene polymer consisting essentially of syndiotactic 1,2-polybutadiene by the successive steps of:\n(a) preparing a catalyst component solution by dissolving, in an inert organic solvent containing 1,3-butadiene, a cobalt compound, soluble in the organic solvent, such as (i) cobalt-.beta.-diketone complex, (ii) cobalt-.beta.-keto acid ester complex, (iii) cobalt salt of organic carboxylic acid, and (iv) halogenated cobalt-ligand compound complex, and an organoaluminum compound,\n(b) preparing a catalyst composition by mixing the catalyst component solution (prepared in step a) with an alcohol, ketone or aldehyde compound and carbon disulfide,\n(c) providing a polymerization mixture containing desired amounts of 1,3-butadiene, the catalyst composition and an inert organic solvent, and\n(d) polymerizing 1,3-butadiene at a temperature which is within the range of -20.degree. C. to 90.degree. C.\nU.S. Pat. No. 3,901,868 indicates that the melting point of the SPBD produced varies in response to the proportion of alcohol, ketone or aldehyde in the polymerization mixture. U.S. Pat. No. 4,153,767 shows that amide compounds, such as N,N-dimethylformamide, can be used in solution polymerizations to reduce the melting point of SPBD being synthesized.\nU.S. Pat. No. 4,429,085 discloses a process for producing syndiotactic 1,2-polybutadiene by suspension polymerization in an aqueous medium. In this aqueous polymerization process polybutadiene which has an essentially syndiotactic 1,2-microstructure is made by the steps of:\n(A) preparing a catalyst component solution by dissolving, in an inert organic solvent containing 1,3-butadiene (a) at least one cobalt compound selected from the group consisting of (i) .beta.-diketone complexes of cobalt, (ii) .beta.-keto acid ester complexes of cobalt, (iii) cobalt salts of organic carboxylic acids having 6 to 15 carbon atoms, and (iv) complexes of halogenated cobalt compounds of the formula CoX.sub.n, wherein X represents a halogen atom and n represents 2 or 3, with an organic compound selected from the group consisting of tertiary amine alcohols, tertiary phosphines, ketones and N,N-dialkylamides, and (b) at least one organoaluminum compound of the formula AlR.sub.3, wherein R represents a hydrocarbon radical of 1 to 6 carbon atoms;\n(B) preparing a reaction mixture by mixing said catalyst component solution with a 1,3-butadiene/water mixture containing desired amounts of said 1,3-butadiene;\n(C) preparing a polymerization mixture by mixing carbon disulfide throughout said reaction mixture, and\n(D) polymerizing said 1,3-butadiene in said polymerization mixture into polybutadiene while agitating said polymerization mixture.\nU.S. Pat. No. 4,751,275 discloses a process for the preparation of SPBD by the solution polymerization of 1,3-butadiene in a hydrocarbon polymerization medium, such as benzene, toluene, cyclohexane, or n-hexane. The catalyst system used in this solution polymerization contains a chromium-III compound which is soluble in hydrocarbons, a trialkylaluminum compound, and a dialkylphosphite, such as di-neopentylphosphite or di-butylphosphite.\nU.S. Pat. No. 4,902,741 and U.S. Pat. No. 5,021,381 discloses a process for preparing a syndiotactic 1,2-polybutadiene latex by emulsion polymerization which comprises polymerizing 1,3-butadiene monomer in an aqueous reaction mixture which is comprised of (1) water, (2) at least one emulsifier, (3) 1,3-butadiene monomer, (4) a catalyst emulsion composition which is prepared by dissolving in an inert organic solvent containing at least one polyene (a) at least one cobalt compound selected from the group consisting of (i) .beta.-ketone complexes of cobalt, (ii) .beta.-keto acid ester complexes of cobalt, (iii) cobalt salts of organic carboxylic acids having 6 to 15 carbon atoms, and (iv) complexes of halogenated cobalt compounds of the formula CoX.sub.n, wherein X represents a halogen atom and n represents 2 or 3, with an organic compound selected from the group consisting of tertiary amine alcohols, tertiary phosphines, ketones and N,N-dialkylamides, and (b) at least one organoaluminum compound of the formula AlR.sub.3 wherein R represents a hydrocarbon radical of 1 to 6 carbon atoms to produce a catalyst component solution, and microfluidizing the catalyst component solution with an oil, a surfactant, and water to an average particle size which is within the range of about 10 nanometers to about 1000 nanometers; and (5) at least one member selected from the group consisting of carbon disulfide and phenyl isothiocyanate.\nThe synthesis of SPBD in an aqueous medium offers several important advantages over solution polymerizations. Water, as a medium in which to carry out such a polymerization, is less expensive, more easily purified, less sensitive to oxygen, and has a higher heat capacity. Conducting such polymerizations in an aqueous medium also permits for higher monomer and higher solids concentrations because of the lower viscosity of a polymer suspension or emulsion compared with that of a polymer solution. The main drawback associated with aqueous suspension and emulsion polymerizations for producing SPBD is the difficulty associated with reducing the melting point of the SPBD. In other words, it is difficult to control the chemical structure and hence the crystallinity and melting point of SPBD which is synthesized in an aqueous medium. Even though numerous modifiers can be used to reduce the level of crystallinity and resulting melting point of SPBD which is synthesized in solution, there are few efficient modifiers for reducing the crystallinity of SPBD which is synthesized in an aqueous medium.\nU.S. Pat. No. 5,011,896 discloses the use of 4-(alkylamino)benzaldehydes, 4-(dialkylamino)benzaldehydes, 2,4-di-(alkoxy)benzaldehydes, 2,6-di-(alkoxy)benzaldehydes, 2,4,6-tri-(alkoxy)benzaldehydes, and 4-(1-azacycloalkyl)benzaldehydes, and 4-(1-azacycloalkyl)benzaldehydes as modifiers for reducing the melting point of SPBD which is synthesized in an aqueous medium. Nevertheless, there is still a need for more highly efficient modifiers which can be used on a commercial basis."}
{"text": "1. Field of the Invention\nThe present invention relates to a monitoring apparatus configured to monitor an image forming apparatus, the image forming apparatus monitored by the monitoring apparatus, and a monitoring system including the image forming apparatus and the monitoring apparatus.\n2. Description of the Related Art\nA known system for monitoring success or failure of change of a network setting and recovering the original setting in case of failure is available.\nJapanese Patent Laid-Open No. 2001-251337 discloses a system in which a device instructed by a server to change its network setting recovers the network setting values before the change in the case where no instruction is given from the server within a predetermined period of time after the change of the network setting.\nThe technique described in Japanese Patent Laid-Open No. 2001-251337 is suitable for the case where the server is controlling the change of the network setting. However, the technique is not suitable for the case where a key role in changing the network setting is played not by the server, but by a customer at which an image forming apparatus is placed.\nIn a system where an image forming apparatus at a customer's site is remotely monitored, various services are provided on the basis of counter information and error information communicated from the image forming apparatus to a remote monitoring server via the Internet. Because a network setting of the image forming apparatus can be implemented easily by the customer on an operation screen of the image forming apparatus, a network breakdown easily occurs due to an incorrect or mistaken change of the setting. In the remote monitoring system, a user of the image forming apparatus does not generally pay attention to the fact that the remote monitoring server and the image forming apparatus are performing network communication with each other. Thus, such an incorrect change of the setting may be left for a while without being corrected. In contrast, since the remote monitoring server detects the occurrence of a network breakdown or the like on the basis of the presence of once-per-day regular communication from the image forming apparatus, there may be a maximum delay of one day from the occurrence to the detection of such a communication breakdown. The remote monitoring system must detect a network breakdown as soon as possible, specify the cause of the network breakdown, and prompts an appropriate person in charge of maintenance of the image forming apparatus to take appropriate measures to correct the network breakdown because of the characteristics of such services as error monitoring and toner inventory management at the image forming apparatus.\nAs in the related art, in case of network breakdown, it is not suitable for the image forming apparatus to recover the original setting or to display an error on a screen of the image forming apparatus because of the following reasons.\nSince a change of the network setting of the image forming apparatus is implemented by a system administrator in change of customer information, changing the setting for the sake of the remote monitoring server can be regarded as ultra vires.\nFurther, in the case where an error is displayed when the image forming apparatus' print operation including a copy operation except for communication with the remote monitoring server is normal, the user may be misled into believing that the image forming apparatus is broken although the operation used by the user is normal. An exemplary network setting that does not influence the use of the image forming apparatus and that hinders communication between the image forming apparatus and the remote monitoring server is a proxy server. A proxy server is used to establish connection from an intranet environment to an Internet environment."}
{"text": "1. Field of the Invention\nThe present invention relates generally to ground transportation management, and in particular, to a method and apparatus for advanced ground transportation management.\n2. Background of the Invention\nMany state and local agencies use Geographical Information System (GIS) databases to manage, plan, and record geographical information in their jurisdictions. For example, the placement of roads, sewers, and other municipal information that are used for planning and management purposes are kept in GIS databases. However, these GIS databases are used only to map these geographical data points for realty purposes, e.g., to know where a public road ends and a private road begins, to know where a sewer line is for purposes of repair, etc. Each municipality typically updates these databases as repairs are undertaken and completed.\nMunicipalities also operate safety departments such as police, fire, and paramedic services. These departments are not provided access to the GIS databases for the associated municipality, and, as such, are unaware of any changes in the database that may affect their operations or assist in managing the operations they control. For example, paramedics may be unaware that a given street is closed for repairs, and be delayed in responding to a call because the paramedics en route to an accident scene tried to use the street that is closed.\nFurther, current routing systems perform routings based upon static speed data. They do not take into account the dynamically changing traffic situation. At best they merely report a status, and are not integrated with a GIS system for use in planning purposes. Many mapping databases report that there is an accident on a given freeway, but do not determine any time of travel on the road, segment, or interval containing the accident. Further, these routing systems are generically determined based on only one data input, namely, a road closure. These systems do not take into account other factors such as equipment status or time of travel between two given points on the roads, segments, alleys, etc. that connect these two points. These systems also do not retain data for analysis after events have occurred to root out systemic problems or determine corrective actions.\nThe large GIS databases, even if combined with other services and data, do not have the capability to provide information to commercial and consumer markets for use in managing fleet and personal travel itineraries. Such access would provide lower fuel costs and shorter travel times, as well as better management of fleet resources.\nEven if the GIS databases were combined with existing services, the number of sensors and other data sources used to augment the GIS databases do not provide proper coverage to accurately predict or determine the optimal route between two points. Even in large metropolitan areas, the percentage of roads monitored by sensors is a small fraction of the number of roads that are in service, and, as such, the data available cannot provide an accurate model of real-time traffic conditions.\nEmergency management operations, typically deployed during times of evacuation, do not utilize GIS databases. Some typical reasons for evacuation, including hurricanes threatening an area, wildfires, biological, nuclear, or chemical attacks, have fixed evacuation routes, and use the same evacuation routes for all different types of emergencies. Emergency operations centers typically do not have access to the tools necessary to dynamically identify the optimal routes for evacuation. As such, there are typically signs marking predetermined roadways as “evacuation routes” rather than dynamic determinations of what route may be best at any given time or for any given emergency. More complex incidents, such as wildfires and terrorist attacks, are more dynamic in nature, and the optimal evacuation plan cannot be predicted due to uncertainties in how the emergency will unfold prior to the actual event.\nFrom the foregoing, it can be seen, then, that there is a need in the art for interconnectivity between the GIS databases and other sources of data. It can also be seen, then, that there is a need in the art to provide access to the combined GIS database for management and operations beyond the municipal schema for use by emergency personnel to determine evacuation routes. It can also be seen that there is a need in the art for a method of dynamically determining evacuation routes based on the imminent or ongoing emergency."}
{"text": "There is significant evidence in the medical literature that repositioning the papillary muscles within the ventricles of the heart during atrioventrical valve repair surgery improves outcomes. The displacement of the papillary muscles, due to ischemia, heart failure, or other causes of ventricular reshaping, results in tethering of the valve leaflets which interferes with their normal functioning. Repairs that focus only on the valve annulus often result in recurrence of regurgitation due to leaflet tethering.\nNumerous methods of papillary muscle repositioning are described in the medical and patent literature. However, these methods of papillary muscle repositioning are typically performed during an open heart surgery. Therefore, there is a need for a less invasive device and method of performing papillary muscle repositioning. In particular, there is a need for a device and method of performing papillary muscle repositioning via a catheter. Moreover, although the literature suggests papillary muscle repositioning, suggested devices and methods have had little commercial success. There is therefore a need for improved devices and methods, regardless of whether delivered via a catheter or in some other way."}
{"text": "The spine includes lumbar vertebra, thoracic vertebrae and cervical vertebrae, protects the spinal cord and nerve root, supports the body and relieves the external impact. The spine of the human body has a structure formed by stacking multiple vertebrae. The intervertebral disc is located between the vertebrae and evenly distributes the load and impact to the entire spine.\nHowever, recently, spinal diseases continue to increase due to the lack of exercise and incorrect postures. The treatment of diseases related to the spine generally includes an indirect treatment using physical therapy and a direct treatment for correcting and fixing the spine by attaching a separate fixing device to the injured vertebrae.\nThere are various causes of spine injury. Among the causes, the bio-mechanic cause is known as the most likely cause. Therefore, for the purpose of the physical chiropractic of the injured spine, a spinal fixing device such as a pedicle screw is generally operated.\nA general spinal fixing operation using the pedicle screw, that is, the above-mentioned direct treatment, is performed by inserting (at least two) pedicle screws into the vertebrae respectively and then by fixing and coupling the housing coupled with the each pedicle screw to the spine rod disposed approximately parallel with the longitudinal direction of the spine.\nMeanwhile, the pedicle screw is divided into a mono type where the screw and the housing are integrally formed with each other and a polyaxial type where the screw and the housing are separately made such that a predetermined rotation is allowed between the screw and the housing, and then are coupled to each other.\nThe shape of the vertebra (into which the pedicle screw is inserted) is not constant according to patients and operation position. Therefore, at the time of the operation, the polyaxial type pedicle screw is widely used in order that a doctor is allowed to insert the pedicle screw more easily to a part into which the mono type pedicle screw cannot be inserted."}
{"text": "1. Field of Invention\nThis invention relates to medical retrieval devices for removing objects from a body, particularly calculi from the urinary and biliary systems.\n2. Description of Prior Art\nMedical instruments are currently in use which reduce the invasiveness and potential trauma previously associated with various medical procedures. One such procedure is the removal of objects, such as kidney stones and gallstones, from the body. Various surgical devices are available which allow objects to be removed from the body without requiring major surgery. One type of surgical device is a mechanical retrieval device. Typically, such instruments consist of 2 or more flexible elements that are joined at their proximal ends and may or may not be joined at their distal ends. The flexible elements are formed in the shape of a basket, cage, grasper, or other entrapping configuration. This basket is attached to a drive wire that passes through the lumen of a small diameter (typically 2.3 mm (7.0 Fr) or less) flexible sheath, which is usually greater than 50 cm in length. The proximal ends of the sheath and drive wire are attached to a multi-part handle, normally constructed of thermoplastic materials, which can typically be operated by the user with a single hand. By manipulating the handle, the drive wire can be pulled back relative to the sheath, collapsing the basket as it retracts into the sheath. In this closed position, the sheath can be passed through the working channel of an endoscope to the proximity of the object to be removed within the patient's body. The basket is expanded to the open position by manipulating the handle, which remains outside the endoscope and the patient's body. The device is then manipulated using the handle until the object becomes enclosed within the basket. This manipulation may include advancing, withdrawing and/or rotating the basket in order to get the object to pass between the flexible elements that comprise the basket. When the object has been successfully engaged within the basket, the endoscope and the retrieval device containing the object are then simultaneously removed from the body.\nA number of designs for the handle of medical retrieval devices are in use. Typically, these handles consist of two main elements, a handle base and an actuation mechanism. The sheath is attached to one of these elements, and the drive wire is attached to the other. By moving the actuation mechanism relative to the handle base, the basket can be retracted into the sheath and extended from it. The handle design that appears to be preferred, based on actual current use and sales volumes, is of a thumb slide design. This design features a main handle base that remains stationary and a thumb slide actuator that slides along a portion of the handle body and has a thumb pad. This handle is held in one hand by wrapping the four fingers of the hand partially around the handle base. The thumb of the same hand is placed on the thumb pad. The device is actuated in one direction by moving the thumb pad away from the proximal end of the handle base, and in the other direction by moving the thumb pad toward the proximal end of the handle base. This type of handle can normally be held in such a way that the range of motion of the thumb required to fully actuate the device is located in a natural and comfortable area near the thumb's resting position. Since considerable skill and dexterity can be required of the user in order to retrieve an object, user comfort is of primary importance.\nThere are many variations of the thumb slide handle in use. However, these and other prior art handle designs do not have a mechanism for rotating the basket to facilitate capturing the object. Rotation can only be accomplished by rotating the entire handle. This method has a number of disadvantages. By rotating the entire handle, the user must accordingly rotate their hand. Since the hand would initially be placed in a natural position, the position of the hand after rotation would not necessarily be comfortable for further manipulation of the handle. Again, user comfort is significant due to the considerable skill and dexterity needed to successfully complete a stone retrieval procedure. Additionally, since the entire handle must be rotated in order to rotate the basket, the sheath must rotate as well as the drive wire. This is a disadvantage because the friction between the sheath and the endoscope's working channel can prevent a smooth 1:1 torque ratio between the handle and the basket. This is particularly relevant when the endoscope is flexible and is in an articulated position. Lack of precise rotational control can increase the difficulty of engaging the object in the basket, thereby lengthening the procedure.\nCertain handle designs have been used which allow rotation of the basket without rotating the entire handle and thus the sheath as well. U.S. Pat. No. 4,046,150 (1977) to Schwartz et al. discloses a retrieval basket with such a handle. This handle has a first member that is attached to the sheath. A second member, which is attached to the drive wire, is located at the proximal end of the first member. The device is actuated by sliding the second member into and out of the first member, which is held stationary. The basket is rotated relative to the sheath by rotating the second member relative to the first member. However, this handle is not of the preferred thumb slide style, and requires the use of two hands to actuate. Also, the actuational and rotational controls are not independent of each other.\nU.S. Pat. No. 5,957,932 (1999) to Bates et al. discloses a retrieval basket with yet another type of handle. This handle is of a pistol grip style, with a control knob located at the proximal end of the handle. The sheath is attached to the main body of the handle, and the drive wire is attached to the control knob. The device is actuated by pulling the control knob out from the main body of the handle and pushing it in. The control knob can also be rotated to rotate the basket. When the trigger portion of the handle is squeezed, mechanical advantage is applied to the actuation of the device. This design has several disadvantages. It requires two hands for normal actuation. And since the control knob is used for both actuation and rotation of the basket, the actuation and rotation are not independent of each other. It also uses a larger number of parts than other handle designs and is therefore is more complex and more expensive to manufacture. The above patents to Schwartz and Bates are incorporated herein by reference.\nAnother type of prior art handle, which is not referenced in any patents, is shown in FIGS. 7 and 8. A handle assembly 190 consists of a stationary portion or handle base 110 and a thumb slide 130 with a thumb pad 131. Handle base 110 has a distal end 181, a proximal end 182, and a length 180. Thumb slide 130 fits partially within handle base 110 and extends out from distal end 181, with thumb pad 131 remaining beyond distal end 181. A hollow tube or sheath 150 has a working length 186 and is attached to thumb slide 130. A drive wire (not shown) passes through the lumen of sheath 150 attaches at the proximal end to handle base 110 and at the distal end to a basket 160. To operate this device, handle base 110 is held in one hand with the four fingers of the hand. The thumb of the same hand is placed on thumb pad 131. When the thumb is extended away from the hand, thumb slide 130 slides out from distal end 181 of handle base 110. This results in the device being in the closed position, as shown in FIG. 7. When the thumb is pulled back toward the hand, thumb pad 131 slides toward distal end 181 of handle base 110. This results in sheath 150 being pulled back to expose basket 160. When thumb pad 131 is pulled back completely to handle base 110, the device is in the open position, as shown in FIG. 8. To rotate basket 160, handle base 110 is rotated relative to thumb slide 130. This design has the disadvantage that the actuation mechanism and the rotation mechanism are not independent. Both actuation and rotation are achieved by movement of the thumb slide and the handle base relative to each other. This design also has the disadvantage that the thumb pad is located beyond the distal end of the stationary handle base. This is a disadvantage because manipulation of the thumb slide is done with the user's thumb extended away from the hand in a somewhat awkward position, which results in less than ideal tactile control over the actuation of the device.\nThe prior art handle designs that do allow rotation of the basket without rotating the entire handle have the disadvantage that the rotation mechanism and actuation mechanism are not independent. As a result, while the basket is being rotated, it may be inadvertently and undesirably expanded or retracted, or while it is being expanded or retracted, it may be inadvertently and undesirably rotated. The retrieval of an object from within a patient's body using an endoscope and a retrieval device is a precise and delicate procedure that requires considerable user skill and dexterity. Since the user's control of the basket is limited by the handle of the device, it is desirable that the handle allows precise and independent control of both the actuation and rotation of the basket, and is comfortable to use."}
{"text": "Virtual machine high availability (referred to herein simply as “high availability,” or HA) and hypervisor-converged object-based (HC/OB) storage are two emerging technologies in the field of computer virtualization. HA is designed to minimize virtual machine (VM) downtime by monitoring the availability of host systems and VMs in a host cluster. If an outage, such as a host or network failure, causes one or more VMs to stop executing, HA detects the outage and automatically restarts the affected VMs on active host systems in the cluster. In this way, HA ensures that guest applications running within the VMs continue to remain operational throughout the outage. One exemplary HA implementation is described in commonly-assigned U.S. Patent Application Publication No. 2012/0278801, published Nov. 1, 2012, entitled “Maintaining High Availability of a Group of Virtual Machines Using Heartbeat Messages.”\nHC/OB storage is a distributed, software-based storage technology that leverages the local or direct attached storage resources (e.g., solid state disks, spinning hard disks, etc.) of host systems in a host cluster by aggregating these locally-attached resources into a single, logical storage pool. Thus, this technology effectively re-purposes the host cluster to also act as a distributed storage cluster. A hypervisor-based storage system layer (referred to herein generically as a “VSAN layer” comprising “VSAN modules”) manages the logical storage pool and enables interactions between the logical storage pool and storage clients, such as VMs running on host systems in the cluster. For example, the VSAN layer allows the VMs to access the logical storage pool during VM runtime in order to store and retrieve persistent VM data (e.g., virtual disk data).\nThe qualifier “object-based” in “hypervisor-converged object-based storage” refers to the manner in which VMs are maintained within HC/OB storage—in particular, the state of each VM is organized as a hierarchical collection of distinct storage objects (or simply “objects”). For example, the files that hold the metadata/configuration of a VM may reside in a file system that is created within a namespace object (also known as a “file system object”), the virtual disks of the VM may reside in virtual disk objects, and so on. Each of these storage objects may be composed of multiple component objects. The VSAN layer provisions, manages, and monitors each of these storage objects individually. For instance, in order to meet a particular storage policy for a particular virtual disk VMDK1, the VSAN layer may determine that the component storage objects that make up the virtual disk object corresponding to VMDK1 should be striped across the locally-attached storage of three different host systems. Through these and other mechanisms, HC/OB storage can provide improved ease of management, scalability, and resource utilization over traditional storage solutions. One exemplary implementation of an HC/OB storage system is described in commonly-assigned U.S. patent application Ser. No. 14/010,293, filed Aug. 26, 2013, entitled “Scalable Distributed Storage Architecture.”\nUnlike non-object-based storage systems, the state of a VM is not contained within a larger, coarse storage container (e.g., a LUN). Having such storage containers provide a couple of benefits. First, a coarse storage container provides a convenient location to store information common to all VMs that use the container. For example, it is possible to create a file system on top of a LUN, create a directory within the file system for each VM whose state is stored on the underlying storage device(s), and then create a directory at the root to store shared information. Second, for a given class of failures, one can reason about the availability/accessibility of all of the VM data stored within a storage container by reasoning about the availability/accessibility of the container itself. For instance, one can determine whether a network failure impacts the accessibility of the VM data by determining if the container is accessible. As a result, there is no need to track the accessibility of each individual VM stored in a single storage container—instead, it is sufficient to track the accessibility of the container itself.\nThe lack of coarse storage containers raises unique challenges when attempting to use HC/OB storage and HA concurrently in the same virtualized compute environment. As one example, existing HA implementations typically maintain information known as “HA protection state” that identifies the VMs in a host cluster that should be failed-over/restarted in the event of a failure. The “master” HA module in the cluster (i.e., the HA module that is responsible for detecting failures and orchestrating VM failovers/restarts) manages this HA protection state by persisting it to a centralized file (or set of files) on the storage tier. If there is an outage that affects a subset of host systems in the cluster, one or more new master HA modules may be elected. Each newly elected master HA module may then retrieve the file from the storage tier to determine which VMs are HA protected. This approach works well if the storage tier is implemented using dedicated shared storage, since the HA protection file can be placed in the storage container storing the configurations for the protected VMs. On the other hand, if the storage tier is implemented using HC/OB storage, there is no convenient location to store such information that is shared across VMs.\nAs another example, in existing HA implementations, when a master HA module detects a failure that requires one or more VMs to be failed-over/restarted, the master HA module executes a conventional failover workflow that involves (1) identifying active host systems for placing the VMs that can meet the VMs' resource needs, and (2) initiating VM restarts on the identified host systems. If the VMs are stored on dedicated shared storage, these two steps are generally sufficient for successfully completing the failover. However, if the VMs are stored on HC/OB storage, there may be cases where a VM cannot be restarted because one or more of its storage objects are not yet accessible to the host system executing on the master HA module (and/or to the host system on which the restart is being attempted). This situation cannot be uncovered using conventional coarse-grained storage accessibility checks. This, in turn, can cause the conventional failover workflow to break down, or result in multiple continuous restart attempts, which can increase the load on the affected host systems.\nAs yet another example, there are certain types of network partitions that can further complicate the HA protection state persistence and VM failover/restart workflows noted above. As one example, if there is a failure that causes the VSAN modules to observe a partition while the HA modules do not, there may be instances where the host system on which the master HA module is running does not have access/visibility to a particular VM (and thus cannot update/retrieve HA protection state information for the VM, or determine its accessibility for failover purposes), while the host systems of other, slave HA modules do have such access/visibility.\nAccordingly, it would be desirable to have techniques for integrating HA with distributed object-based storage systems like HC/OB storage that overcome these, and other similar, issues."}
{"text": "According to some technical analysts, there will be over 50 billion connected “things” by the year 2020. This will completely transform current infrastructures and will drive new innovations in industry, products, and services. Internet-of-Things (IoT) is term that represents devices and systems that communicate over a network, such as the internet. The devices and systems may include sensors.\nExposure to certain gasses is detrimental to human health. Depending on the level and duration of exposure, significant health effects may be incurred, up to and including death. Exposure to these gasses occurs when individuals are in proximity to sources of the gas. Sensors are used to measure gas concentrations."}
{"text": "Fibrous batts, mats or boards are used to thermally insulate various surfaces, such as the inside or outside of pipes and ducts, as well as refrigerators, air conditioners, furnaces, automobile hoods and the like. Heating, cooling or ventilating systems usually use air ducts, through which flows air as the heating and cooling medium. These ducts are constructed in various shapes, for example as round tubes or having a rectangular shape. They are usually made of sheet metals, such as galvanized steel, or plastics, or fiberglass reinforced plastics, or cellulosic materials such as wood or fiberboard, and the like. The construction and method of making these ducts and conduits is described in U.S. Pat. Nos. 3,092,529; 3,212,529; and 3,394,737. The inside of these ducts is usually black so as not to see the duct through the vents. The innner surface of these ducts should be smooth, to allow the air to flow with as little friction as possible, to save energy in transporting the air. This is usually achieved by first coating the fiberglass insulation which is then installed in the air ducts. The coating of the fiberglass insulation normally does not interfere with the excellent noise dampening qualities of the fiberglass insulation. The duct noise is often generated by the fan and by the flowing air, or, in the case of refrigerators, automobiles, furnaces and the like, by the engines and motors. Generally, the fiberglass mats, or batts are coated with an aqueous coating composition, often after the surface of the fiberglass mat has first been covered with a woven or nonwoven fabric.\nConventionally, the coatings are applied by spraying the coating composition onto the surface of the fiberglass batting. The coating is then dried in heated ovens, or by radiant heat. These procedures are generally carried out continuously. The finished fiberglass batting is then pulled up and often compressed by applying vacuum before shipping in order to save on shipping costs.\nOne such coating composition is disclosed in U.S. Pat. No. 3,926,894. The coating composition of U.S. Pat. No. 3,926,894 consists of a halogenated organophosphorous plasticizer, a latex binder and a mineral filler. However, a serious disadvantage of such a coating composition is the relative volatility of the plasticizer, especially in air ducts where large volumes of heated air pass over the coating, which makes the coating environmentally undesirable because the vapors emanating from the plasticizers are quite toxic.\nAccordingly, it is an object of this invention to provide a coating having good mechanical properties, low volatility of the ingredients, fire retardation, as well as low black smoke generation. Furthermore, it is an object of this invention to apply the coating more efficiently to the fiberglass insulation by whipping air into the coating composition and applying the coating as a heat collapsible foam by roller or knife coating."}
{"text": "The present disclosure is generally related to methods of forming metal oxide nanostructures and nanostructures thereof, in particular, titania nanostructures.\nTitania is a well-know material with a broad range of applications including photonic crystals, photocatalysts, and photovoltaic cells. While several methods are known for nanostructuring titania, including thermal imprinting, many challenges remain mainly due to the properties of commonly-used sol-gel type titania precursors. The sol-gel type titania precursors are formed at low pH (approximately 1), are generally highly reactive and moisture sensitive, and form gels. They are often diluted in organic solvents during the sol-gel reaction to mitigate gelation, which causes large volume shrinkages during the nanostructuring process. Further, they are usually highly viscous and require high pressure for the nanostructuring process. Prior work directed toward thermal imprinting titania nanostructures used sol-gel type precursors [C. Goh, K. M. Coakley, M. D. McGehee, Nano Lett. 5, 1545 (2005), P. Yang, T. Deng, D. Zhao, P. Feng, D. Pine, B. F. Chmelka, G. M. Whitesides, G. D. Stucky, Science, 282, 2244, (1998)] or a mixture of titania colloidal particles and polymers [M. Wang, H.-G. Braun, and E. Meyer, Chem. Mater. 14, 4812 (2002)]."}
{"text": "1. Field of the Invention\nThe invention relates to a surface-mount light-emitting diode (LED). More particularly, the invention relates to a surface-mounted LED including an LED chip mounted in a recess provided in a substrate, in which a sealant resin is applied to seal the LED chip.\n2. Description of the Related Art\nA conventional surface-mounted LED may have the structure shown in FIGS. 8 and 9. FIG. 8 is a top view and FIG. 9 is a cross-sectional view taken along line A-A in FIG. 8. It includes a pair of circuit patterns 51a, 51b formed in the upper surface of an insulator 50 at both opposite ends. The circuit patterns 51a, 51b extend from edges of the insulator 50 to the lower surface around sides of the insulator 50. A recess 52 is provided almost at the center in the upper surface of the insulator 50. A circuit pattern is formed over the entire bottom 53 and the entire inner circumferential surface 54 in the recess 52. This circuit pattern is connected as an extension from the one circuit pattern 51a of the pair of circuit patterns 51a, 51b formed in the upper surface of the insulator 50. The other circuit pattern 51b extends toward the center of the insulator 50.\nAn LED chip 56 is mounted on the bottom 53 in the recess 52 via a conductive adhesive 55. In this case, the lower electrode of the LED chip 56 is connected to the circuit pattern 51a formed on the bottom 53 in the recess 52 for achievement of the electrical conduction therebetween. The upper electrode of the LED chip 56 is connected via a bonding wire 57 to the circuit pattern 51b that extends substantially toward the center of the substrate for achievement of the electrical conduction therebetween.\nA light transmissive resin 58 is applied to cover the LED chip 56 and the bonding wire 57 for sealing and protecting the LED chip 56 from the external environmental elements such as humidity, dirt, and gases. In addition, it protects the bonding wire 57 from mechanical stresses such as vibrations and impacts (see, for example, Patent Publication 1: JP-A 7/202271 along with its associated English Abstract, machine translation, and drawings, which are submitted herewith in an Information Disclosure Statement and are hereby incorporated in their entireties by reference).\nSurface-mounted LEDs are often employed together with other surface-mounted electronic components and are generally surface-mounted on a component-mounting board in an electronic instrument through a solder reflow furnace. In this case, the surface-mounted LED is extremely small and, accordingly, the temperature of the whole LED may elevate almost up to the heating temperature of the solder reflow furnace.\nAt the elevated temperature, the light transmissive resin for use in sealing the LED chip and the bonding wire have a difference in thermal expansion coefficient as compared to the circuit pattern formed on the bottom in the recess. This difference and other factors may cause a stress that results in peel at a contact interface between the circuit pattern and the resin. Then, the light transmissive resin exerts a force on the conductive adhesive and the LED chip to lift them above the circuit pattern. In this case, the light transmissive resin may peel off the circuit pattern, possibly resulting in an electrical property failure.\nEven after the surface-mounted LED is mounted on the component-mounting board, the repetition of switching the LED on/off can cause repeated thermal expansion and contraction of the light transmissive resin. In addition, the stress at that time may exert the same action as described above and possibly result in an electrical property failure."}
{"text": "1. Field of the invention\nThe present invention relates to a rear wheel suspension controller for a vehicle, particularly to a rear wheel suspension controller which is effective against a sporadic shock caused by a protrusion or a sinking of a road surface on which an automobile is running.\n2. Prior art\nConventionally, the spring constant, damping force, bush characteristic or stabilizer characteristic of each of various suspension components provided between a body of a vehicle and its wheels is altered under control depending on conditions of a road surface or running conditions of the vehicle in order to prevent the vehicle from being shocked or vibrated and keep the controllability and the stability of the vehicle good. For example, altering the spring constant of the air spring of a suspension depending on conditions of the road surface, altering the damping force of a shock absorber, and simply making the characteristic of a bush or a stabilizer variable were proposed in published unexamined Japanese patent applications No. sho 59-23712 and No. sho 59-26638, in those No. sho 58-30542 and No. sho 59-23712, and in Japanese utility model application No. sho 58-26605 and published unexamined Japanese utility model application No. sho 59-129613, respectively. In such control, when it is detected by a vehicle height sensor that the vehicle is running on a rough road or when it is detected by a brake sensor or an accelerator sensor that the front of the vehicle has gone up or down, the characteristic of each suspension of the vehicle is altered to maintain a good controllability and stability of the vehicle running on the rough road, or to prevent the front of the vehicle from going up or down further. However, under the above-mentioned conventional control, the vehicle is not judged to be running on a rough road, until a large turbulence is continuously detected by the vehicle height sensor. When the vehicle is judged to be running on a rough road, the spring constants of the suspensions for all the wheels of the vehicle or the damping forces of the shock absorbers for all the wheels are increased to produce a desired effect. If the vehicle passes over a joint of road patches or its sporadic protrusion or sinking, the vehicle usually receives only one shock and resumes running on a flat part of the road again, so that the characteristic of each suspension is not altered. For that reason, passengers of the vehicle are not protected from an unpleasant shock due to such sporadic protrusion or sinking, which is different from the case that the vehicle is running on a rough road having continuous protrusions or sinkings. In some cases of passing over such sporadic protrusion or sinking, the controllability and the stability of the vehicle deteriorate as well."}
{"text": "Augmented reality is a technology in which a person's conception of reality can be enhanced, typically through augmented sound, video or graphics displays. The augmentation is typically implemented via various technologies, such as a headset that may be worn by the person. One or more augmented views may be presented to the person through the headset.\nThe augmented reality headset typically includes a wearable computer and an optical display mounted to the headset. The wearable computer may include a wireless telecommunication capability, permitting a wireless connection from the wearable computer to a server computer. Because of the wireless telecommunication capability, the augmented reality headset may be used to interact with the server computer to accomplish various tasks."}
{"text": "1. Field of the Invention\nThe present invention relates broadly to management of medical information. More specifically, the present invention relates to management of medical information to perform and report measurements of physician efficiency.\n2. The Prior Art\nRecent evidence has suggested that about 10-20% of physicians, across specialty types, practice inefficiently. Efficient means using the appropriate amount of medical resources to treat a medical condition and achieve a desired health outcome. Thus, efficiency is a function of unit price, volume of service, intensity of service, and quality of service. This group of inefficient physicians is responsible for driving 10% to 20% of the unnecessary, excess medical expenditures incurred by employers and other healthcare purchasers, equating to billions of dollars nationally.\nTo improve market efficiency, it is useful to apply a system that accurately measures individual physician efficiency. Recent evidence has demonstrated that leading physician efficiency measurement systems have only about 15-30% agreement across measurement systems. This means that when one system ranks a physician as inefficient, only about 15-30% of the other systems ranked the same physician as inefficient. The remaining 70% (or more) of systems ranked the same physician as efficient.\nThese findings show that existing systems have significant error in attempting to accurately identify inefficient physicians. The error needs to be eliminated, or significantly reduced, if healthcare purchasers are to accurately identify inefficient physicians and take action (e.g., attempt to change physician behavior, provide incentives for employees to use more efficient physicians). Every physician falsely measured as efficient (or inefficient) leads to continued inefficiency in the healthcare marketplace.\nThere are ten common physician (or physician group) efficiency measurement errors present in most existing physician efficiency measurement systems, which are in order of importance: (1) examine all episodes of care for a physician; (2) use a physician's actual episode composition; (3) no severity-of-illness measure by medical condition; (4) no identification of different episode treatment stages; (5) no age category assignment by medical condition; (6) no tracking mechanism for related complication episodes; (7) improper episode outlier criteria; (8) under-report charges attributed to partial episodes; (8) over-report charges attributed to episode endpoints; and (10) no minimum number of episodes of care. These errors are discussed next.\nMany physician efficiency methodologies continue to examine “services per 1,000 members” or “all episodes of care” tracked to a physician. These approaches probably add the most to efficiency measurement error. The methodologies attempt to adjust services per 1,000 members and to adjust all episodes of care by age and gender—and then compare one physician's utilization pattern to a peer group average. However, age and gender explain less than 5% of the variance in a patient's medical expenditures. This means that over 95% of the variance is unexplained, and may be attributed to differences in patient health status.\nSome methodologies adjust services per 1,000 members and adjust all episodes of care based on specific International Clinical Modification of Diseases ninth edition (ICD.9) code algorithms that measure expected resource intensity. The idea is that a patient's diagnosis codes will provide more predictive power than age and gender alone. The most predictive of the published and marketed models explain only 20% to 30% of the variance in a patient's medical expenditures. This means that 70% or more of the variance continues to be unexplained, and may be attributed to differences in patient health status.\nPhysicians often criticize the services per 1,000 members and the all episodes of care methodologies that use a predictive case-mix adjustment factor. Physicians state that the methodologies do not appropriately adjust for differences in patient health status—rightly stating that their patients may be “sicker.”\nIf all claim line items (CLIs) or episodes of care tracked to a physician are used in the efficiency analysis, then up to 70% of the observed utilization difference between physicians may be attributed to patient health status differences. Therefore, patient health status differences are measured rather than individual physician efficiency differences. This weakness in current case-mix adjustment tools means that not all CLIs or patient episodes of care treated by a physician can be examined. Instead, an isolated set of more prevalent medical conditions by severity-of-illness level needs to be examined across physicians of a similar specialty type.\nThe second measurement error, which occurs in most if not all current efficiency measurement systems, occurs when the physician's actual episode composition is used. The reason is as follows. The differences in physicians' patient case-mix composition results in differences in variability (i.e., the standard deviation) around a physician's average episode treatment charges. This variability is not due to the efficiency or inefficiency of a physician, but instead results because longer and more resource-intensive medical conditions generally require more services and, therefore, have more potential variability around average (or mean) episode treatment charges.\nFor example, easier-to-treat upper respiratory infection (URI) episodes may have the following mean and standard deviation (with outlier episodes removed): $185±$65. Here, the standard deviation around the mean is not large—and is 0.35 the size of the mean (i.e., 65/185=0.35). However, easier-to-treat pediatric asthma episodes may have the following mean and standard deviation (with outlier episodes removed): $1,650±$850. Here, the standard deviation around the mean is larger than for URI episodes—and is 0.52 the size of the mean (i.e., 850/1,650=0.52).\nThe variation difference between the two conditions is 49% greater for asthma than URIs_[(0.52−0.35)/0.35]. This variation difference occurs for two reasons: (1) more resource-intensive conditions require more services to treat; and (2) there generally are a small number of episodes available to examine in a given physician efficiency study as compared to the universe of episodes that could actually be studied—and a smaller number of episodes results in a higher chance for variability around the mean. This variation is not the result of physician treatment pattern differences.\nIf the statistically based variability around the mean is not corrected, then substantial error may enter into the physician efficiency measurement equation. Consequently, the final physician efficiency score differences may be attributed to the statistical condition-specific variability around the mean episode charge (due to the case-mix of episodes treated).\nThe above example showed that the variation difference may be 50% or more (around a condition-specific mean episode value). Logically, then, if all episodes treated by physicians are examined and efficiency scores are calculated, there has to be some statistical bias present.\nA significant statistical bias may be present. Using a more traditional episode-based efficiency measurement methodology, lower-episode-volume physicians treating patients with a higher case-mix index score are more likely to receive an inefficient ranking as compared to lower-episode-volume physicians treating patients with a lower case-mix index score. This finding results because a physician with higher case-mix patients treats episodes having more variability (i.e., a greater standard deviation) around average episode treatment charges. With a low volume of episodes (most often the norm, and not the exception), this physician needs only a few higher-cost episodes then the peer group average to make his/her treatment pattern appear significantly higher than the peer group comparator.\nHowever, a physician with lower case-mix patients treats episodes having less variability around average episode treatment charges. With a low volume of episodes, this physician's treatment pattern will not be as influenced by one or two higher-cost episodes as compared to the peer group average. Consequently, his/her treatment pattern does not appear (as often) significantly higher than the peer group comparator.\nThus, by examining all medical condition episodes, a substantial component of any observed physician efficiency difference may be attributed to statistical condition-specific variability around the mean episode charge—and not to physician treatment patterns efficiency. This effect may be present even when we examine the easier-to-treat episodes (SOI-1 level episodes) for the medical conditions.\nThe third error takes place in those efficiency measurement systems that do not employ an appropriate episode severity-of-illness measure. Severity-of-illness may be defined as the probability of loss of function due to a specific medical condition. Most, if not all, current claims-based episode groupers and methods do not have an appropriate severity-of-illness index by medical condition. Consequently, significant clinical heterogeneity remains in many episodes for a given medical condition. The end result may be physician efficiency differences that are attributed to inaccurate episode severity-of-illness adjustment, and not to physician treatment patterns variation.\nMoreover, some claims-based episode groupers stratify formulated episodes for a medical condition by the presence or absence of a specific surgery or service (e.g., knee derangement with and without surgery; ischemic heart disease with and without heart catheterization). The reason for performing this stratification is to reduce episode heterogeneity for a medical condition. In effect, the stratification serves as a sort of severity-of-illness adjustment.\nHowever, stratification based on the presence of surgery or a high-cost service results in at least two physician efficiency measurement errors: (1) performing surgery versus not performing surgery is the treatment patterns variation we need to examine in determining physician efficiency, and this variation is not captured in more traditional methodologies; and (2) the episodes of care are unnecessarily separated into smaller groups whereby physicians may not have enough episodes to examine in any one smaller group. Consequently, the stratified episodes of care need to be recombined for accurate physician efficiency measurement.\nThe fourth physician efficiency measurement error occurs in claims-based episode groupers do not have a method for identifying different episode treatment stages including initial, active, and follow-up treatment stages. Identifying different treatment stages is important in medical conditions, such as breast cancer, prostate cancer, colorectal cancer, acute myocardial infarction, and lymphoma. For example, breast cancer should be stratified into initial, active, and follow-up treatment stages.\nAn initial breast cancer episode is one where the patient has a surgery for the cancer (e.g., lumpectomy, modified radial mastectomy). An active breast cancer episode is one where no surgery is present, but chemotherapy or radiation treatment is observed within the episode. Here, the patient underwent surgery in a previous study period, so no surgical event is found in the patient's current ongoing breast cancer episode. Instead, during the study period, the claims data shows that the patient is being treated with chemotherapy and/or radiation. The presence of these treatments defines an active breast cancer episode. The utilization pattern and charges are different for an active breast cancer patient as compared to an initial breast cancer patient. A follow-up breast cancer episode is one where no surgery, chemotherapy, or radiation treatment is present in the patient's episode of care. After initial and active treatments, physicians will continue to code for breast cancer over the future years of patient follow-up care.\nIn a given study period, physicians do not treat an equal distribution of each episode type (initial, active, and follow-up). Moreover, the episode types have different average charges. About 20% of episodes may be classified as initial breast cancer episodes. Overall care for initial breast cancer episodes ranges between $15,000 and $25,000 per episode. About 15% of episodes may be classified as active breast cancer episodes. Overall care for active breast cancer episodes ranges between $12,000 and $18,000 per episode. About 65% of episodes may be classified as follow-up breast cancer episodes. Overall care for follow-up breast cancer episodes ranges between $350 and $600 per episode.\nConsequently, the blending of the three treatment stage episodes results in average treatment charges of about $5,500 to $6,500 per episode. In fact, this is the average breast cancer charge that would be observed for most claims-based episode groupers.\nThe blending of initial, active, and follow-up episodes may lead to substantial physician efficiency measurement error. For example, assume during a study period that Oncologist A treats mostly active breast cancer patients, while some other oncologists have a good mixture of active and follow-up patients. Then, Oncologist A's treatment pattern for breast cancer will appear inefficient (as compared to his peer group of oncologists) because active episodes are about 30 times more expensive to treat than follow-up episodes. In fact, Oncologist A's treatment pattern difference is attributed to a different treatment stage episode case-mix.\nTherefore, treatment stage episode types need to be correctly identified and separately examined. Otherwise, the final physician efficiency score differences may be attributed to nothing more than the initial, active, and follow-up episode case-mix.\nThe fifth error happens in those physician efficiency measurement systems that do not examine condition-specific episodes by age category. Studies have illustrated that broad-based age bands are important to separately examine—even after episodes have been assigned a severity-of-illness index. The reason is that physicians tend to treat children and adults differently for most conditions. For example, children are less likely than adults to receive a chest x-ray and potent antibiotics for many medical conditions. If episodes are not examined by broad-based age category, the end result may be physician efficiency differences that are attributed to patient age differences—and not to treatment patterns variation.\nThe sixth error occurs in those physician efficiency measurement systems that do not link and include the charges and utilization from a patient's complication episodes to his underlying medical condition. Complications are those episodes that are clinically related to the underlying medical condition. Consequently, many condition-specific episodes have under-reported charges. In fact, actual outputs from some claims-based episode groupers may show under-reported charges for patients with diabetes and other chronic conditions (e.g., asthma, congestive heart failure).\nFor example, the reason for the under-reported episode charges is that physicians code up to 70% of an average diabetic's charges under related complications to the diabetes (e.g., eye, neuropathies, circulatory, renal) and not diabetes care. Therefore, without considering and including related complication episodes with the actual diabetes episode, physician efficiency differences may be attributed to incomplete episode charges and utilization—and not to treatment pattern variations.\nFurthermore, for patients with specific medical conditions, any model that attempts to stratify patients by health risk may produce unstable or erroneous results. The reason is that a patient is missing key claims information needed to accurately classify a patient into an appropriate severity-of-illness and other classes. For example, without tracking related complications to a diabetic patient, many diabetic patients will appear to have no complications when in fact they have eye or circulatory complications.\nThe seventh physician efficiency measurement error happens when the condition-specific outlier episode analysis is not performed in an appropriate manner. Many current methodologies perform the high-end outlier analysis by eliminating a percent of condition-specific episodes at the peer group (or aggregate episode) level. That is, the methodologies eliminate the high-end outliers before assigning episodes to physicians.\nHowever, this method results in physician efficiency measurement error because a higher proportion of episodes assigned to the most inefficient physicians will be eliminated (as compared to the proportion of episodes eliminated for efficient physicians). Consequently, the inefficient physicians' condition-specific treatment patterns now more closely resemble the treatment patterns of the efficient physicians.\nAn example demonstrates this error. Assume Physician A has seven episodes of acute bronchitis with the following per episode charges: $235, $245, $325 $400, $525, $550, and $600. Also, the outlier cut-off threshold for high-end outlier episodes is set at $399 at the peer group level. Physician A now has only three episodes remaining at $235, $245, $325. The mean charge is $268 per episode. Assume Physician B also has 7 episodes of acute bronchitis with the following per episode charges: $210, $225, $235, $255, $285, $320, and $390. The peer-group level outlier threshold remains at $399. Therefore, Physician B has all seven episodes remaining, and the mean charge is $274.\nThe end result shows no statistical difference between Physicians A and B. The mean episode charge of Physician A is slightly lower than Physician B (i.e., $268 versus $274 per episode). However, using an outlier rule where we eliminate 5% of episodes (or at least 1 high-end outlier) are eliminated at the physician level, the results are significantly different. Physician A now has six remaining episodes (i.e., here we eliminate only 1 high-end outlier), and the mean charge of the six non-outlier episodes is now $380 per episode. For Physician B, the mean charge for the six non-outlier episodes is now $255 per episode. Physician A is statistically higher in average (or mean) episode charges than Physician B by $125 per episode.\nThe eighth error occurs in those systems that under-report charges attributed to partial (or incomplete) episodes of care. Some methodologies do not separate partial from complete episodes of care when measuring physician efficiency. Partial episodes result because a patient enrolled in a health plan during the study period or disenrolled during the study period. However, including partial episodes leads to inaccurate efficiency measurement because of under-reported episode charges—especially when some physicians have more partial episodes than other physicians.\nA reason partial episodes often slip through the cracks and into an efficiency analysis is because the methodologies do not use a membership eligibility file to ensure the member is present for the entire study period. Instead, the methods assume that a condition-specific episode of care is complete if the episode exceeds some minimum duration time period. For example, if a patient's episode of diabetes is 40 days or more in duration, then the episode is marked as complete—and not partial. If a patient's diabetes episode is 39 days or less, then the episode is marked as partial.\nApplying an indiscriminate time period duration to condition-specific episodes produces a high percentage of episodes marked as complete, which are actually partial (or incomplete) episodes. That is, many health plan's have membership turnover rates of 20% or higher. Consequently, a diabetes episode of 40 days duration—marked as complete—has at least a 20% chance of being a partial episode of care because of membership turnover. The end result may be physician efficiency differences that are attributed to the inclusion of partial episodes—and not to treatment patterns variation.\nThe ninth error happens in physician efficiency measurement systems that over-report charges attributed to episode endpoints. Some methodologies do not appropriately end a patient's episode of care before measuring a physician's efficiency. For example, chronic conditions may continue indefinitely and, therefore, patient episodes of care may be of various durations (e.g., 60 days or 600 days)—depending on the amount of available patient claims data. The end result may be physician efficiency differences that are attributed to excessively long or variable chronic condition episode durations—and not to treatment patterns variation.\nThe tenth error takes place in those systems that impose few requirements for having a minimum number of episodes in a certain number of medical conditions. Many methodologies do not require a minimum number of condition-specific episodes when comparing a physician's efficiency to a peer group. Instead, only a small handful (e.g., less than 10 episodes) are enough. However, there may be significant episode of care heterogeneity in one or two condition-specific episodes—even after applying a sophisticated severity-of-illness index. Consequently, examining an episode here-and-there for a physician may introduce significant error into a physician's efficiency measurement. The end result may be physician efficiency differences that are attributed to the heterogeneity in the low number of episodes examined—and not to treatment patterns variation.\nVarious systems have been patented in the episode of care field. Such systems are shown, for example, in U.S. Pat. Nos. 5,557,514, 5,835,897 and 5,970,463. However, none of these systems adequately overcome the aforementioned problems with respect to appropriately building and analyzing episodes of care. As importantly, existing systems fail to discuss an episode-of-care-based system for measuring individual or physician group efficiency measurement."}
{"text": "An example of a scaleable computing solution is a partitionable computing system. In such a system a number of elements (e.g., computing cells) can be combined into a partition that is dedicated to perform a specific computing function. Multiple partitions can exist in the same partitionable computing system, each having a specific function. A malicious attack on one partition could result in the entire partitionable system being compromised."}
{"text": "1. Field of the Invention\nThis invention relates to an improved field effect transistor circuit and more particularly to an enhancement/depletion mode FET circuit for providing an output voltage swing suitable for turning off depletion mode FET devices."}
{"text": "The liquefaction of low boiling point gases, such as air and the components of air, such as oxygen, nitrogen and argon, has been practiced for over 100 years, and the liquefaction of such gases on an industrial scale has been practiced since the beginning of the 20th century. Typically, commercial liquefiers are designed to produce hundreds of tons of liquid cryogens per day. Such industrial liquefiers are reliable, and are capable of producing liquefied gas with relatively high energy efficiency. For consumers of liquefied gas requiring relatively small quantities, small insulated containers, known as dewars, are filled with liquefied gas produced by commercial facilities and transported to the consumer. Consumers of small quantities of liquefied gas include hospitals, which require oxygen for delivery to patients and nitrogen for use as a refrigerant. Also, people suffering from chronic respiratory insufficiency that have been prescribed home oxygen by their physicians may have liquefied oxygen delivered to their residences.\nInitially, attempts to provide such a liquefier involved efforts to miniaturize large scale liquefying plants. However, due to the complexity of such systems, which are typically based on the Claude cycle or its variants, these attempts failed. Also, the extremely small mechanical components resulting from the miniaturization of such liquefiers were expensive to produce and unreliable in operation. Current liquefiers often involve complex and/or expensive liquefaction components, and often lack safety features to make a liquefaction system safer for residential, small-scale, and/or portable use.\nFor the above-stated reasons, it would be advantageous to provide a method and apparatus for improving the safety, efficiency, and/or cost of producing and storing relatively small quantities of liquefied gas at the location where the liquefied gas is to be used, such as at an oxygen therapy patient's residence."}
{"text": "At present, tetrahydrofuran is a widely used cyclic ether. The primary purposes of tetrahydrofuran are as solvent and binder of various resins, as solvent and extraction solvent of printing inks, as the surfactant of synthetic leather and as the material for synthesizing elastic fibers of polytetrahydrofuran. The purpose of other cyclic ethers has not been fully developed yet.\nThe process of synthesizing cyclic ethers from alkanediols and the catalysts used therein have been well described in the known literatures and patent references. For example, those wherein phosphoric acid is used as a catalyst are disclosed in U.S. Pat. Nos. 2,251,292, 2,251,835 and 4,124,600; those wherein surfuric acid is used as a catalyst are disclosed in Ger. Offen. 2,509,968, U.S. Pat. Nos. 4,665,205 and 5,099,039, Jpn. Tokkyo Koho 78-43,505 and 78-43,506; those wherein aluminum oxide is used as a catalyst are disclosed in U.S. Pat. No. 4,196,130, Brit. 508,548, Ger. Offen. 2,856,455 and USSR SU 1,158,562.\nThe process of synthesizing cyclic ethers from alkanediols and the catalysts used therein can be classified into two types: one is a process of homogeneous reaction, the other is a process of non-homogeneous reaction. In the process of homogenous reaction, sulfuric acid and phosphoric acid are the representative catalysts. The drawback of such a process is that part of the catalyst is distilled out with the reaction product, resulting in difficulties of separating the catalyst from the reaction product. Moreover, the acidities of sulfuric acid and phosphoric acid catalysts are strong enough to severely corrode the reactor. In addition, when sulfuric acid and phosphoric acid are used as the catalyst, side reactions easily result in which a great amount of coke is produced. The existence of the coke in the reacting liquid may affect the activity of the catalyst. Therefore, in the continuous synthesis process, while the reacting liquid should be removed from the reactor and part of the reacting liquid should be treated, the problem of acidic waste also results. As to the process of non-homogeneous reaction, aluminum oxide is the representative catalyst. The drawback of this process is that the aluminum oxide catalyst is able to exhibit strong activity at high temperatures only, usually at temperatures higher than 250.degree. C. Such high temperatures generally far exceed the boiling points of many alkanediols. Therefore, the reaction must be carried out in a gasous state and the cost of the equipment and operation will be increased to no avail.\nThe conversion of alkanediols into cyclic ethers using LaHY, CaHY and H-ZSM-6 zeolites was proposed by C. P. Bezouhanova and F. A. Jabur et al. in React. Kinet. Catal. Lett., Vol. 51, No. 1, pp. 177-181 (1993). However, the reaction was carried out in the condition of a gaseous state at temperatures higher than the boiling points of the reactants. The poor selectivity of cyclic ethers is a drawback of this process. It is not satisfactory because the further step of purification is necessary in the subsequent work-up.\nIn view of these drawbacks of the above-mentioned conventional arts, the inventor has studied intensively and found that they can be resolved by using a crystalline aluminosilicate zeolite as the catalyst which allows the reaction to be carried out in the liquid phase. Therefore, this invention is able to be achieved."}
{"text": "It is common for companies to sponsor sporting events or to otherwise place their advertisements within a sports arena. For example, a company may contract with a party having rights in an arena, team or league to place a banner within a stadium during game days, to place a logo on a team jersey, to have an advertisement displayed on digital signage within a stadium, etc. Sponsors and holders of rights in the advertising space often determine pricing and desirability of specific advertising space based in part on in-person audience attendance at the sporting event and a size of the television audience watching the sporting event at home. However, it is increasingly common, due in part to changes in the way that people consume content, that these attendance and television viewership numbers may significantly underestimate the number of people that actually saw at least a clip or highlight of the sporting event that contained a sponsor's logo or advertisement. For example, short video highlights are often played across many different television channels as well as shared on the Internet via social media networks, video sharing platforms and other services. These additional exposures are not typically tracked in any reliable or comprehensive manner."}
{"text": "1. Field of the Invention\nThe present invention relates to a clasp utilizing the attraction of a permanent magnet, and more particularly to obviating such inconveniences as disruption of magnetic records by magnetic flux from such clasps by preventing the flux lines of the magnet from leaking externally.\n2. Description of the Prior Art\nThere have been proposed various clasps which utilize the attraction of a permanent magnet, and almost all of them are conventionally aimed at effectively using the magnetic attraction of a permanent magnet rather than taking counter-measures to prevent disruption of magnetic records by the permanent magnet.\nAs great innovations have been achieved in the recording technology recently and magnetic recording means such as tapes, cards or notes have become household items, safeguarding those means demands special attention and care.\nWe have entered an era where articles having magnetic records are used daily, such as various magnetic tapes and magnetic disks to tickets for transportation, admission tickets, or cash cards for bank accounts.\nThe content of such magnetic records as recording tapes, etc., however, can be lost because they are easily destroyed when placed under the influence of the magnetic flux lines of a magnet and the occurrence of such a disruption can not easily be observed from outside appearance."}
{"text": "Characteristic of cardiothoracic surgery is the post-operative patient who is sent to the Intensive Care Unit (ICU) intubated due to respiratory requirements. Approximately half of these patients are extubated within their first twenty-four post-operative hours. In most cases these patients are extubated within the first three days. There are some, however, who remain intubated for a significant length of time. When a surgeon identifies a patient who will require intubation longer than seven days, the surgeon will usually decide to perform a tracheotomy on that patient. The breathing support tube enters the trachea rather than entering the mouth for the trached patient. Communication for a intubated or trached patient is minimal due to the inability to speak resulting in the patient, hospital staff and loved ones resorting to the reading of lips, nodding of heads and squeezing of hands to communicate.\nWithout effective communication, the intubated or trached patient may not receive the standard of care he or she would otherwise receive had he or she been able to effectively communicate. The lack of communication also creates unnecessary levels of anxiety which the patient must endure. Nurses and hospital staff ask many questions from the patient pertaining to their prognosis and progress which may never get fully or even adequately answered. A doctor or nurse is not able to treat a symptom which they know little or nothing about. In addition, other problems arise due to the insufficient communication from the patient. Localized areas of pain are often mis-diagnosed, resulting in over-medication generally or the medication of an area which is not the source of pain. Proper and essential treatment given in an adequate and timely manner will help resolve or prevent many post-operative complications and decrease the patient\"\"s length of stay in the hospital. This begins with providing the patient a clear and precise means of communication.\nAccordingly, there has been a need for an ICU communication device which in the immediate post-operative period can provide assistance to an intubated or trached patient. What is also needed is a device which provides the communicating elements necessary over the patient\"\"s post-operative stay in the hospital with not only with medical care providers but also with visiting family and loved ones. Further, a communication device is needed which accomplishes the desired function while being easy to manufacture and use while remaining cost effective. The present invention fulfills these needs and provides other related advantages.\nThe present invention resides in a device which facilitates communication between a voice-disabled patient and his or her care provider and others. The device comprises, generally, a housing having at least one display surface, indicia displayable on the display surface, and a marker associated with the housing. The indicia may be utilized by the patient to indicate the status and needs of the patient. The marker is usable by the patient to communicate to a third party the patient\"\"s status and needs utilizing the indicia.\nIn one form of the invention, the device for communicating with a voice-disabled patient comprises a clipboard having at least one eraser-board surface, and an erasable marker attachable to the clipboard, and indicia imprinted onto the eraser-board surface. More particularly, the clipboard includes two eraser-board surfaces, and an eraser is connectable to the clipboard.\nIn another form of the invention, an electronic device for communicating with a voice-disabled patient comprises a housing, a computer within the housing, a touch pad visual screen disposed on the housing and in electronic communication with the computer, and computerized screen layouts generated by the computer having touch activated icons indicating the patient\"\"s status and needs. The housing includes handles, and a speaker is disposed within the housing for audibly transmitting a computerized voice corresponding to the icons displayed on the visual screen.\nIn both embodiments the indicia includes descriptive words and phrases, and graphical representations of a human body. The indicia may further include a grid containing alphabetical letters, numbers and universal symbols, and a pain scale. Moreover, the words and phrases may include the patient\"\"s physical and emotional status, and the graphical representations of the human body may have correlating descriptive words and phrases indicating the physical status of specific parts of the body.\nOther features and advantages of the present invention will become apparent from the following more detailed description, taken in conjunction with the accompanying drawings which illustrate, by way of example, the principles of the invention."}
{"text": "Contact sensors have been used to gather information concerning contact or near-contact between two surfaces in medical applications, such as dentistry, podiatry, and in the development of prostheses, as well as in industrial applications, for instance to determine load and uniformity of pressure between mating surfaces, and in the development of bearings and gaskets. In general, these sensors include pressure-sensitive films designed to be placed between mating surfaces. These film sensors, while generally suitable for examining static contact characteristics between two generally flat surfaces, have presented many difficulties in other situations. For instance, when examining contact data between more complex surfaces, surfaces including complex curvatures, for instance, it can be difficult to conform the films to fit the surfaces without degrading the sensor's performance.\nMore serious problems exist with these materials as well. For instance, film-based contact sensor devices and methods introduce a foreign material having some thickness between the mating surfaces, which can change the contact characteristic of the junction and overestimate the contact areas between the two surfaces. Moreover, the ability to examine real time, dynamic contact characteristics is practically non-existent with these types of sensors.\nA better understanding of the contact conditions at joints and junctions could lead to reduced wear in materials, better fit between mating surfaces, and longer life expectancy for machined parts. For example, one of the leading causes of failure in total joint replacement prostheses is due to loosening of the implant induced by wear debris particles worn from the polymeric bearing component. A better understanding of the contact conditions between the joint components would lead to reduced implant wear and longer implant life.\nWhat is needed in the art are contact sensors that can provide more accurate and/or dynamic contact information concerning a junction formed between two surfaces of any surface shape."}
{"text": "This invention relates to a resistive film made from a paste comprising a powdery, vitreous carbon as an electrically conductive component in admixture with an electrically non-conductive component. This invention also relates to a method for producing a vitreous carbon in powder form by pyrolyzing a resin in inert atmosphere.\nVitreous carbon has a Mohs' hardness of 6. Therefore, reducing it to a powder entails significant effort.\nProduction of powdery, vitreous carbon is disclosed in German patent document DE 27 18 308 A1, in which acrylamides mixed with water-soluble salts are pyrolyzed, and a vitreous carbon component is recovered following pyrolysis by dissolving a salt component in water. The vitreous carbon is then dried, and a powder produced in this manner is further pulverized as needed. The method is very time-consuming.\nUse of an precursor polymer having a three-dimensional, cross-linked structure for producing a vitreous carbon is disclosed in Plastverarbeiter, vol. 41, no. 6, pages 16-21 (1990). After shaping via casting or molding, the polymer is cured and additionally machined down. No mention is made of possible pulverization of the vitreous carbon.\nGerman patent document DE 30 02 112 A1 discloses a paste for producing polymer-film integrated circuits having predetermined electrical conductivity. A predetermined electrical resistivity of the polymer-film integrated circuits to be produced is achieved by mixing electrically conductive and electrically non-conductive film components. A desired resistivity value is thus produced by adding the electrically non-conductive component.\nIt is generally known that the electrical conductivity of a paste, or of a film made of a paste, formed as a mixture of an electrically conductive component and an electrically non-conductive component, is largely determined by the specific electrical conductivity and the concentration, i.e. a packing density, of the conductive component in the film system. Electrical resistivity in the film increases nearly exponentially when the conductive component reaches a critical minimum concentration in the film. When the electrically conductive component is increased, electrical resistivity stabilizes once an optimal concentration is achieved. When a specific resistivity value is established through proportional increases in the electrically non-conductive component, electrical and mechanical stability may diminish as a result, owing to lack of homogeneity of the film.\nWhen this paste is used for a resistive film, on a potentiometer, the electrically conductive component, for example carbon, causes a typical electrical micro-heterogeneity of the surface of the film, and thus causes increased contact resistivity at the film's wearing surface.\nEuropean patent document EP 0 399 295 A1 discloses a use of a vitreous carbon as an electrically conductive component in a resistive film. An electrical resistivity level is established by modifying a concentration or packing density of an electrically conductive component in the resistive film with respect to that of an electrically non-conductive component, such as a binding agent. In order to maintain desirable properties, such as mechanical and electrical stability, an optimum concentration or packing density of electrically conductive particles of the vitreous carbon, i.e. the mixture ratio of the two components, can vary only within a limited range.\nIt is an object of this invention to provide a method for producing powdery, vitreous carbon that saves money and time. A further object of this invention is to furnish a paste having predetermined electrical conductivity, and to disclose a resistive film having a high degree of abrasion resistance as well as mechanical and electrical stability."}
{"text": "The optimum mercury vapor pressure for production of 2537 .ANG. radiation to excite a phosphor coating in a fluorescent lamp is approximately six millitorr, corresponding to a mercury reservoir temperature of approximately 40.degree. C. Conventional tubular fluorescent lamps operate at a power density (i.e., typically measured as power input per phosphor area) and in a fixture configuration to ensure operation of the lamp at or about a mercury vapor pressure of six millitorr (typically in a range from approximately four to seven millitorr); that is, the lamp and fixture are designed such that the coolest location (i.e., cold spot) of the fluorescent lamp is approximately 40.degree. C. Compact fluorescent lamps, however, including electrodeless solenoidal electric field (SEF) fluorescent discharge lamps, operate at higher power densities with a cold spot temperature typically exceeding 50.degree. C. As a result, the mercury vapor pressure is higher than the optimum four to seven millitorr range, and the luminous output of the lamp is decreased.\nOne approach to controlling the mercury vapor pressure in an SEF lamp is to use an alloy capable of absorbing mercury from its gaseous phase in varying amounts, depending upon temperature. Alloys capable of forming amalgams with mercury have been found to be particularly useful. The mercury vapor pressure of such an amalgam at a given temperature is lower than the mercury vapor pressure of pure liquid mercury.\nUnfortunately, accurate placement and retention of an amalgam to achieve a mercury vapor pressure in the optimum range in an SEF lamp are difficult. For stable long-term operation, the amalgam should be placed and retained in a relatively cool location with minimal temperature variation.\nCommonly assigned U.S. Pat. No. 4,262,231 of Anderson et al., issued Apr. 14, 1981, which is incorporated by reference herein, describes situating a lead-tin-bismuth amalgam in an electrodeless SEF fluorescent lamp by wetting the amalgam to a metal wire structure, such as a helical structure or a cylindrical screen, which is fixed within the tip-off region of a lamp envelope. Alternatively, Anderson et al. describe melting the amalgam onto an indium-coated, phosphor-free portion of the interior surface of the lamp envelope.\nSmeelen U.S. Pat. No. 4,622,495 describes another scheme for locating an amalgam within an electrodeless SEF fluorescent lamp by attaching an amalgam holder to a tubular indentation (hereinafter referred to as a re-entrant cavity) within the lamp envelope. Disadvantageously, this requires a glass-to-metal seal; and a reliable glass-to-metal seal is difficult to achieve in manufacturing.\nAccordingly, it is desirable to provide a relatively simple method for locating an amalgam in an electrodeless SEF fluorescent discharge lamp which provides an optimal operating location for the amalgam, while not requiring a glass-to-metal seal or an internal amalgam holder. Moreover, the amalgam should be held in place during lamp manufacturing without significantly interfering with other lamp processing steps."}
{"text": "The invention relates to piezo-electric actuating elements for recording heads in mosaic-type recording equipment, in particular those wherein a recording ink or fluid contained in a compression chamber, surrounded by the cylindrical actuating element, is forced out in drops through piezo-electric constriction of the actuating element.\nThe use of the piezo effect to operate recording heads has been known for a long time. German Inspection specification No. 2,405,584, for example, describes a pulse drop injection system in which a glass tube is encircled by a piezo-electric transducer which constricts in synchronism with a pulse generator and thereby forces recording fluid in the glass tube to be discharged drop by drop.\nRecording tubes of this type are supported in holding devices within which the entire electro-mechanical transducer is firmly embraced by clamping means, with the external electrode of the transducer being secured to such holding device. Consequently, a loosening of the recording tube and holding device, as a result of any change in the external diameter resulting from constriction of the transducer, must be avoided as such loosening would automatically become noticeable from the resulting deterioration in the quality of the recording."}
{"text": "The invention concerns a coated cutting tool. More specifically, the invention pertains to a coated cutting tool wherein the coating scheme for the cutting tool includes a layer of a Group IVB metal-aluminum alloy applied by physical vapor deposition (PVD) and a layer of alumina applied by PVD to the layer of the Group IVB-aluminum alloy.\nAlumina as a bulk material exhibits high oxidation resistance, high chemical stability, high hardness and good wear resistance. Alumina thus has desirable properties for use in material removal, e.g., metal removal, applications. Notwithstanding these desirable properties, the use of alumina cutting tools has certain drawbacks due to the low toughness and poor formability of alumina. To improve the cutting performance, some manufacturers have included additives such as chromium oxide, magnesium oxide, titanium oxide, nickel oxide and refractory metal carbides to the alumina-based cutting tool. The presence of these additives has resulted in alumina cutting tools with improved performance properties.\nAnother way to overcome the disadvantages of alumina and yet still achieve the beneficial properties thereof is to coat tougher and more readily formed cutting tool substrates with an alumina coating via physical vapor deposition (PVD) techniques. The end result of this process is an alumina-coated cutting tool.\nWhile alumina-coated cutting tools achieve certain performance levels, there have been drawbacks with these tools. In this regard, the following documents concern the application of an amorphous alumina coating to a substrate.\nThe article by Knotek et al. entitled \"Sputtered hard materials based on titanium and aluminum for wear protection\", International Pulvermet. Tagung., DDR, Dresden (1985), pp. 181-196, mentions the PVD (magnetron sputtering) application of separate coatings of alumina, TiN and TiAlN. The article states that these coatings have equivalent, or in the case of TiAlN better, performance than CVD coatings.\nThe Sumomogi et al article entitled \"Adhesion Evaluation of RF-Sputtered Aluminum Oxide and Titanium Carbide Thick Films Grown on Carbide Tools\", Thin Solid Films, 79 (1981), pp. 91-100, discusses the PVD application of alumina coatings on cemented carbide tools (ISO P20).\nThe Shinzato et al. article entitled \"Internal Stress in Sputter-Deposited Al.sub.2 O.sub.3 Films\", Thin Solid Films, 97 (1982), pp. 333-337, mentions the application of an alumina coating to cemented carbide (ISO P20), high speed steel, and Corning Pyrex glass substrates via either a conventional r.f. apparatus or a planar magnetron r.f. apparatus.\nThe Kazama et al. article entitled \"Alumina Films Prepared by Ion Plating\" alumina on a WC cutting tool!, mentions the application of an alumina coating via ion plating to a WC substrate to from an alpha-alumina coating.\nJapanese Patent Application S54-2982 to Murayama et al. addresses the use of the ion plating method to apply an alumina (or alumina and TiC) coating to a substrate. Practical Example II concerns the alumina coating on a cutting tool (ISO P30 alloy).\nOne drawback of alumina coated cutting tools pertains to the lack of adequate adhesion between the substrate and the alumina coating or layer. One way to improve the adhesion of the alumina coating is to provide an intermediate layer or layers. The following documents discuss the use of an intermediate layer.\nIn Ramos et al., \"Adhesion improvement of RF-sputtered alumina coatings as determined by the scratch test\", Jour. Adhesion Sci. Technol., Vol. 7, No. 8, pp. 801-811 (1993), the authors looked at the use of a titanium or a titanium nitride intermediate layer between a high speed steel substrate (AISI M2) and an alumina coating applied via RF-sputtering. The overall conclusion was that the presence of a titanium or titanium nitride intermediate layer improves the adhesion of the alumina coating to the high speed steel substrate.\nJapanese Patent Publication No. 57-120667 to Doi et al. mentions in the context of a cutting tool the PVD application of an intermediate layer of alumina over a base layer of TiC. The outer layer is TiN. This patent also mentions a six-layer coating scheme including alumina as the outer layer with TiC, alumina, TiCO as the intermediate layers using a plasma CVD device. Example No. 3 of this patent comprises an eight layer scheme in which the outer layer is alumina and the intermediate layers include TiC, alumina, TiCO and TiCN.\nJapanese Patent Application No. H4-308075 to Matsuda et al. discusses the application by ion plating of an alumina coating to a substrate. The alumina may be applied directly to the surface of the substrate. The alumina may also be applied to an intermediate layer previously formed on the substrate. The intermediate layer may be titanium, titanium carbide or titanium nitride."}
{"text": "The subject matter disclosed herein generally relates to a shot sleeve for a die casting process and, more particularly, to low modulus shot sleeves for high temperature die casting.\nA die casting process utilizes a mold cavity defined between mold parts. Molten metal material is feed in to the mold cavity and held under pressure until the metal hardens. The mold parts are then separated and the cast part removed. In some processes a shot sleeve is utilized to hold molten material and introduce that material to the cavity. The shot sleeve includes an opening for introducing molten material into a bore of the shot sleeve that leads to the mold cavity. A plunger or piston moves within the bore of the shot sleeve to push the molten material through the shot sleeve and inject the molten material into the mold cavity. The piston is subsequently withdrawn and additional material can be introduced into the bore for fabricating another part within the same mold cavity, i.e., the shot sleeve is reused for multiple molding operations (e.g., die casting operations).\nThe shot sleeve can experience very high temperatures due to the molten metal material that is passed through the bore of the shot sleeve. Accordingly, the shot sleeve and/or components thereof are fabricated of materials compatible with such high temperatures. However, materials that are compatible with the high temperatures encountered during the die casting process can be costly and difficult to machine. Further, materials that are compatible with the high temperatures may result in shot sleeves with relatively low life cycles. That is, the high temperatures can lead to failure of the shot sleeves, even when the shot sleeve is formed from high temperature materials. Accordingly, it is desirable to design and develop shot sleeves that can withstand the high temperatures while reducing cost, easing manufacturing, and/or increasing the life cycle of shot sleeves."}
{"text": "The invention relates to method of electrodeposition employing resinous reaction products containing cationic groups. More particularly, this invention relates to cationic resinous reaction products which are prepared from chain extended epoxide resins.\nCationic electrodepositable resins are known in the art. A preferred class of resins are those prepared from epoxy resins such as disclosed in U.S. Pat. No. 4,104,147 to Marchetti, Jerabek and Zwack.\nThis patent discloses chain extension of polyepoxides with organic polyols such as polymeric polyols. The chain extended products can then be reacted with a secondary amine and solubilized with acid to form cationic electrodepositable compositions. The resins have excellent properties such as high rupture voltage, good film forming properties and deposit as films with good flexibility.\nOne problem associated with polymeric polyol chain extension is competing reactions. Under chain extension reaction conditions, which are usually in the presence of an amine catalyst, epoxy-epoxy reaction and epoxy-secondary hydroxyl reactions compete with the desired polymeric hydroxyl-epoxy reaction. The competing reactions may consume too much of the epoxy functionality resulting in the presence of excess amine in the reaction product which adversely affects the dispersion properties of the resin as well as its throw power and film-forming properties. Also, these competing reactions if not controlled can present manufacturing difficulties, for example, undesirably high resin viscosities which are believed to be due to polymer branching.\nAn indication of this can be seen in FIG. 2 which is a plot of the reduced Gardner-Holdt viscosity (50 percent resin solids in 2-ethoxyethanol) versus time in hours of the reaction mixture which involves chain extension of a polyglycidyl ether of a polyphenol with a poly(oxytetramethylene) glycol having a molecular weight of 650. As shown in FIG. 2, the viscosity increases rapidly with time. If the viscosity is not carefully monitored, the reaction mixture could easily go to gelation. This is a possibility in a commercial production situation where the operator in charge of the reaction may not be able to monitor carefully the viscosity of the reaction with time.\nIt has been found that these problems can be significantly minimized by chain extending with a polymercapto compound, particularly a polymeric polymercapto compound. It is believed that under chain extension reaction conditions, the mercapto-epoxy reaction goes in relatively high yield with a minimum amount of competing side reactions. This results in a resin which has better properties and which is easier to manufacture, particularly on a commercial scale.\nReference is made to FIG. 1 which is a plot of the reduced Gardner-Holdt viscosity versus time for Example I of the present invention. Example I involves chain extension of a polyglycidyl ether of a polyphenol with a dimercapto polymer obtained from reacting poly(oxytetramethylene) glycol having a molecular weight of 650 with mercapto propionic acid in a molar ratio of 1:2. As shown in FIG. 1, the viscosity increases relatively slowly over the period of about 11/2 hours. At this point, the viscosity remains essentially constant with time. In commercial production, this can be important because if the operator in charge of the reaction is distracted and loses track of the time of the reaction, the viscosity of the reaction mixture will level out and not proceed to gelation."}
{"text": "This invention relates to a quotation marks or \"comb teeth\" pattern signal generating circuit of a video camera, and in particular to a zebra signal generating circuit of a video camera in which, an electronic view finder (which is abbreviated as E.V.F) is utilized and especially when taking photographs at a field location, the gain band of the video signal is divided into stages the user is able to see and by selecting the appropriate zebra pattern among these stages, an improved video image quality can be obtained.\nIn general, the zebra pattern means a comb teeth pattern, and this pattern is made to couple to the video signal by dividing the frequency down to a quarter of the 3.58 MHZ carrier which is the sub-carrier of the color signal, since the sub-carrier signal is inverted at every 1 Hz the comb teeth pattern appears in the picture screen. Accordingly, the visual status of the object becomes possible to observe in the electronic view finder of the video camera.\nTherefore, when pressing the white balance of a video camera, the zebra signal appears and is produced up to the third stage from the white in the gray scale chart.\nThe zebra signal is conventionally produced when the video signal is a light quantity of over 100 IRE, the zebra signal can indicate the state and range of incident light quantity with regard to the surrounding environment during picture taking and is useful to check the state of the camera, and is used for making a pertinent video signal and photograph by controlling the lens iris in response to the zebra signal. However, the video camera using a conventional zebra pattern is structured to merely warn about the state or light quantity of the picture screen by inputting the zebra pattern when above a certain specific value i.e., only the portion of the light quantity being directed, and being designed not to be able to observe the state and light quantity of the picture screen when below the specific value, and therefore, when taking a picture of the field location with a video camera, because the state and scope of the light quantity is not possible to observe without a measuring instrument or master monitor, there is a problem that a picture of good quality is hard to obtain with only the electronic view finder scanned with black and white."}
{"text": "Drive-in facilities by means of which bank customers may transact their business without leaving their vehicles have become increasingly popular, particularly in suburban locations, to the point where two, three, or even more stations are required to accommodate the traffic; and where multiple stations are involved, it is extremely difficult to provide two or more teller stations in locations where the motorist can drive up to the teller station and transact business with the teller by means of an extensible drawer in which currency, deposit slips and the like may be transferred between the teller and the customer. To alleviate this problem, as well as conserve space and enhance traffic flow, remote customer stations have been provided in the form of islands having remote facilities for transferring items back and forth between the remote customer station and the teller station. In addition to increasing the number of customers who may transact their business at any given time, a lesser number of tellers is usually required in that one teller can service more than one customer station.\nFor the most part, vacuum systems are utilized to convey the items back and forth between the customer stations and the teller stations. Such vacuum systems are expensive to install and to operate, and the number and size of items which can be transported is relatively limited due to the size of the normally cylindrical containers required for travel through the vacuum lines. In addition, should one of the lines be plugged, either by one of the containers or by foreign materials inadvertently, or even deliberately, introduced into the vacuum lines, considerable difficulty is often encountered in removing the obstruction and any containers which might be trapped in the system by reason of the obstruction. Similar objections are encountered with other forms of moving conveyors, particularly those which move underground or are otherwise relatively unaccessible for maintenance and repair; and if a stoppage results due to a malfunction in only a portion of the system, the entire system may have to be shut down until the necessary parts can be obtained and the repairs made, which may put the system out of operation for a number of hours, days, or even weeks.\nIn contrast to the foregoing the instant invention provides an integrated conveyor system which is relatively inexpensive both insofar as initial cost is concerned as well as in its cost of operation, and which can be easily installed and service."}
{"text": "Exercise machines having alternating reciprocating foot supports configured to traverse or travel about a closed path to simulate a striding, running, walking, and/or a climbing motion for the individual using the machine are well known in the art, and are commonly referred to as elliptical exercise machines or elliptical cross-trainers. In general, an elliptical or elliptical-type exercise machine comprises a pair of reciprocating foot supports designed to receive and support the feet of a user. Each reciprocating foot support has at least one end supported for rotational motion about a pivot point or pivot axis, with the other end supported in a manner configured to cause the reciprocating foot support to travel or traverse a closed path, such as a reciprocating elliptical or oblong path or other similar geometric outline. Therefore, upon operation of the exercise machine to rotate the proximal end, each reciprocating foot support is caused” to travel or traverse the closed path. The reciprocating foot supports are configured to be out of phase with one another by 180° in order to simulate a proper and natural alternating stride motion.\nAn individual may utilize an elliptical or elliptical-type exercise machine by placing his or her feet onto the reciprocating foot supports. The individual may then actuate the exercise machine for any desired length of time to cause the reciprocating foot supports to repeatedly travel their respective closed paths, which action effectively results in a series of strides achieved by the individual to obtain exercise, with a low-impact advantage. An elliptical or elliptical-type machine may further comprise mechanisms or systems for increasing the resistance of the motion, and/or for varying the vertical elevation or height of the closed path. In addition, the reciprocating motion of the feet to achieve a series of strides may be complemented by a reciprocating movement of the arms, whether assisted by the exercise machine via a suitably configured mechanism or system, or unassisted.\nA typical closed path may comprise a generally horizontal outline having a longitudinal axis therethrough. Depending upon the exercise machine, a closed path may be many different sizes. As such, a particular measurement of interest to individuals with respect to an elliptical or elliptical-type exercise machine is “stride length”. A stride length is essentially a measurement of the distance separating the two furthest points along the longitudinal axis of the closed path. Therefore, upon actuation of the exercise machine, a single stride may be referred to as travel by the reciprocating foot support, and therefore the foot of a user, along the closed path from a first endpoint on the along the longitudinal axis of the closed path to the a distal endpoint, also on the longitudinal axis. The stride and resulting stride length provided by an exercise machine, although simulated and possibly modified, is comparable to a single stride achieved during natural and/or modified gait of an individual.\nObviously, the strides, and particularly the stride lengths, between different individuals may vary, perhaps considerably. Indeed, a person of small stature will most likely have a much shorter stride length than a person of large stature, and thus will be more comfortable on an exercise machine configured to accommodate his or her particular size and resulting stride length. As such, it is important that the exercise machine function with a stride that corresponds to the stride of the user. The challenge arises when the exercise machine is intended for use by many individuals that may or may not have the same stride length. Moreover, it may be desirable within an exercise routine to vary the speed or frequency of strides along the closed path, the resistance felt, and/or the vertical height of the closed path, wherein some or all of these variable elements may require the user to adapt his or her stride to the changing routine to realize a more natural motion.\nDespite their many advantages, and despite recent efforts to attain such, elliptical or elliptical-type exercise machines are devoid of a simple and efficient way to vary their stride length for the purpose of accommodating the stride lengths of individuals of different size and of providing a more natural stride motion. Many prior related exercise machines exist in the art that comprise complex or intricate solutions. However, many of these are difficult to operate at best, and are also expensive to manufacture and cumbersome to assemble as many of them comprise several components or linkages to ultimately achieve a variable stride length.\nAnother inherent deficiency with the many prior related exercise machines comprising a mechanism or system for varying the stride length of the machine is that they are so complex in design that it would be difficult to utilize the system or mechanism technology on different machines without requiring significant modifications to the machine, if possible at all."}
{"text": "Water softeners find wide applications throughout society. In many applications, it is desirable to soften the water by removing the hardness minerals from the water before use. This is particularly critical in boiler operation where use of hard water will create boiler scale and rapidly reduce operating efficiencies.\nA common water softening process is to use water softeners designed for this purpose. Water softening tanks contain cation exchange resin capable of exchanging hardness ions, i.e., calcium and magnesium for sodium ions which are very soluble.\nWhen the hardness exchanging capacity of the water softening resin has exhausted it stops producing soft water. It then becomes necessary to regenerate the resin with a saturated solution of sodium or potassium chloride. Because of cost, sodium chloride is usually the chemical of choice.\nSodium chloride brine solution is made in a separate tank built and designed for this purpose, and this tank is called a brine tank.\nModern water softeners are well engineered and designed to produce soft water with all regeneration actions done automatically, including the transfer of the saturated brine from the brine tank to the water softener tanks.\nIn order for the water softener resin to be properly rejuvenated, the saturated brine solution must be of high quality and a measured volume must be delivered whenever needed.\nA properly designed and engineered brine tank will provide these needs by delivering a measured quantity of saturated salt brine containing a fixed amount of dissolved salt per gallon of water.\nThis is accomplished by using a horizontal salt grid in a vertical tank. The height and diameter of the salt grid varies for each softening system, depending on many factors, but in all cases the height of the salt grid sets the volume of water in the brine tank.\nIn actual practice, the brine system is set to fill the brine tank with fresh water from the bottom of the tank to approximately 1\" above the salt grid and then shut off.\nUsing this method, only 1\" of water touches the vertical salt pile, which may be several hundred pounds in weight, stored on top of the salt grid.\nThis system is called a dry salt shelf system, as opposed to a wet salt brine tank system where most or all the salt is immersed in water. The dry salt shelf system has significant advantages over the wet salt system. The dry salt shelf method produces 100% saturated brine (specific gravity 1.2) all the time where wet salt methods do not. The dry salt shelf system affects more dry salt storage in the same size brine tank than a wet salt system. A dry salt shelf system is easier to keep clean than the wet salt system. A dry salt shelf system does not require a gravel support bed at the bottom of the brine tank. The dry salt shelf system offers lower maintenance costs to the operator, no gravel cleaning or replacement.\nThe dry salt shelf system has no messy brine float valves as used above the liquid brine on wet salt systems. These float valves become corroded with salt creep and require repair and/or replacement frequently. The dry salt shelf system uses brine float or refill valves in the lower section of the brine tank (below the shelf) and are less exposed to the risk of malfunctions or corrosion, thus operating more efficiently. The dry salt shelf system uses all of the salt stored before the brine tank needs to be refilled. Liquid below the shelf is saturated brine even if only one grain of salt remains on the shelf. The brine tank salt refill is less often with the dry salt shelf system because of the greater salt storage capacity it offers. Brine tank corrosion is reduced or eliminated on steel brine tanks with the dry salt shelf system because the liquid level is down below the dry salt, thus less air/brine exposure. The dry salt shelf system allows more programmed salt delivery scheduling because the salt stored is easily seen and thus the quantity remaining can be easily determined. The dry salt shelf system allows the use of all grades of salt, even the most economical rock type salt. The dry salt shelf type brine system can be cleaned in less than one hour, regardless of size whereas a wet salt tank may take one day and require the water softener to be down.\nDissolving of salt starts immediately and continues until the volume of water beneath the salt grid becomes saturated with dissolved salt. When saturation occurs, dissolving ceases. Stored salt above the salt grid not in contact with the water remains dry, preventing bridging and mushing.\nUsing a salt grid enables an engineer to calculate the quantity and quality of a particular size brine tank will produce. The engineer then is able to select the proper brine tank for the water softener system. It is imperative that the grid and support system be strong to support the mass of weight placed upon it. Until now, salt grids and support systems have usually been made from pegboard. It is readily available and cheap; however, in contact with the salt brine it tends to deform, warp and those portions of the salt grid left unsupported tend to break and collapse, dumping the salt stored on them down into the brine measuring area.\nWhen this occurs, it causes the water softener to malfunction.\nThe salt grid and support system must be rebuilt, and in time it fails again.\nFor these reasons, the salt grid and support system could not be used in larger brine tanks.\nAs the demand increased for larger and larger water softeners, the demand for more saturated brine increased.\nIn order to provide this requirement, brine tanks increased in size and the pegboard salt grid and support system could not be used as they were not strong enough to hold the weight.\nThese larger systems were forced to use the less efficient older method of wet storage.\nThis method consists of loading the brine tank with several hundred pounds of gravel on the bottom. Several hundred pounds of salt is then poured upon the gravel and then adding water until a portion or all of the salt is submerged.\nThe measuring advantage of the salt grid system is lost.\nThe salt brine produced by the wet salt storage method is often of poor quality and submerged salt tends to bridge and mush, causing maintenance problems.\nThe salt grid plate and support system we have designed and built is strong enough that it may be used in the larger systems.\nAgain, the engineer can calculate the exact quantity of brine needed by utilizing the salt grid method."}
{"text": "The present invention comprises a new Ageratum, botanically known as Ageratum houstonanium, and hereinafter referred to by the variety name ‘Agros’.\n‘Agros’ is a product of a planned breeding program. The new cultivar ‘Agros’ has purple flowers, early and continuously flowering and has a habit that is compact, upright while mounded and is freely branching.\n‘Agros’ originated from a hybridization in a controlled breeding program in Enkhuizen, Netherlands. The female parent was an unpatented seedling identified as ‘X52-2’ with light purple color. ‘X52-2’ has a more vigorous habit than ‘Agros’.\nThe male parent of ‘Agros’ was an unpatented hybrid seedling identified as ‘54-1’ with purple color. ‘54-1’ has a more vigorous habit than ‘Agros’.\n‘Agros’ was selected as one flowering plant within the progeny of the stated cross in August 2003. The pollination took place in September 2002 and the seed sowing in March 2003, all in a controlled environment in Enkhuizen, Netherlands.\nThe first act of asexual reproduction of ‘Agros’ was accomplished when vegetative cuttings were taken from the initial selection in August 2003 in a controlled environment in Enkhuizen, Netherlands.\nHorticultural examination of plants grown from cuttings of the plant initiated in early spring 2004 in Enkhuizen, Netherlands; Gilroy, Calif. USA; and Angers, France, and continuing thereafter, has demonstrated that the combination of characteristics as herein disclosed for ‘Agros’ are firmly fixed and are retained through successive generations of asexual reproduction.\n‘Agros’ has not been observed under all possible environmental conditions. The phenotype may vary significantly with variations in environment such as temperature, light intensity and day length.\nA Plant Breeder's Right for this cultivar was granted in the European Union on Oct. 22, 2007. ‘Agros’ has not been made publicly available more than one year prior to the filing of this application."}
{"text": "Many medical procedures involve the extraction and replacement of flowing blood from, and back into, a donor or patient. The reasons for doing this vary, but generally, they involve subjecting the blood to some process that cannot be carried out inside the body. When the blood is outside the patient it is conducted through machinery that processes the blood. The various processes include, but are not limited to, hemodialysis, hemofiltration, hemodiafiltration, blood and blood component collection, plasmaphresis, aphresis, and blood oxygenation.\nOne technique for extracorporeal blood processing employs a single “access,” for example a single needle in the vein of the patient or a fistula. A volume of blood is cyclically drawn through the access at one time, processed, and then returned through the same access at another time. Single access systems are uncommon because they limit the rate of processing to half the capacity permitted by the access. As a result, two-access systems, in which blood is drawn from a first access, called an arterial access, and returned through a second access, called a venous access, are much faster and more common. These accesses include catheters, catheters with subcutaneous ports, fistulas, and grafts.\nThe processes listed above, and others, often involve the movement of large amounts of blood at a very high rate. For example, 500 ml. of blood may be drawn out and replaced every minute, which is about 5% of the patient's entire supply. If a leak occurs in such a system, the patient could be drained of enough blood in a few minutes to cause loss of consciousness with death following soon thereafter. As a result, such extracorporeal blood circuits are normally used in very safe environments, such as hospitals and treatment centers, and attended by highly trained technicians and doctors nearby. Even with close supervision, a number of deaths occur in the United States every year due to undue blood loss from leaks.\nLeaks present a very real risk. Leaks can occur for various reasons, among them: extraction of a needle, disconnection of a luer, poor manufacture of components, cuts in tubing, and leaks in a catheter. However, in terms of current technology, the most reliable solution to this risk, that of direct and constant trained supervision in a safe environment, has an enormous negative impact on the lifestyles of patients who require frequent treatment and on labor requirements of the institutions performing such therapies. Thus, there is a perennial need in the art for ultra-safe systems that can be used in a non-clinical setting and/or without the need for highly trained and expensive staff. Currently, there is great interest in ways of providing systems for patients to use at home. One of the risks for such systems is the danger of leaks. As a result, a number of companies have dedicated resources to the solution of the problem of leak detection.\nIn single-access systems, loss of blood through the patient access and blood circuit can be indirectly detected by detecting the infiltration of air during the draw cycle. Air is typically detected using an ultrasonic air detector on the tubing line, which detects air bubbles in the blood. The detection of air bubbles triggers the system to halt the pump and clamp the line to prevent air bubbles from being injected into the patient. Examples of such systems are described in U.S. Pat. Nos. 3,985,134, 4,614,590, and 5,120,303.\nWhile detection of air infiltration is a reliable technique for detecting leaks in single access systems, the more attractive two-access systems, in which blood is drawn continuously from one access and returned continuously through another, present problems. While a disconnection or leak in the draw line can be sensed by detecting air infiltration, just as with the single needle system, a leak in the return line cannot be so detected. This problem has been addressed in a number of different ways, some of which are generally accepted in the industry.\nThe first level of protection against return line blood loss is the use of locking luers on all connections, as described in International Standard ISO 594-2 which help to minimize the possibility of spontaneous disconnection during treatment. Care in the connection and taping of lines to the patient's bodies is also a known strategy for minimizing this risk.\nA higher level of protection is the provision of venous pressure monitoring, which detects a precipitous decrease in the venous line pressure. This technique is outlined in International Standard IEC 60601-2-16. This approach, although providing some additional protection, is not very robust, because most of the pressure loss in the venous line is in the needle used to access the patient. There is very little pressure change in the venous return line that can be detected in the event of a disconnection, so long as the needle remains attached to the return line. Thus, the pressure signal is very weak. The signal is no stronger for small leaks in the return line, where the pressure changes are too small to be detected with any reliability. One way to compensate for the low pressure signal is to make the system more sensitive, as described in U.S. Pat. No. 6,221,040, but this strategy can cause many false positives. It is inevitable that the sensitivity of the system will have to be traded against the burden of monitoring false alarms. Inevitably this leads to compromises in safety. In addition, pressure sensing methods cannot be used at all for detecting small leaks.\nYet another approach, described for example in PCT application US98/19266, is to place fluid detectors near the patient's access and/or on the floor under the patient. The system responds only after blood has leaked and collected in the vicinity of a fluid detector. A misplaced detector can defeat such a system and the path of a leak cannot be reliably predicted. For instance, a rivulet of blood may adhere to the patient's body and transfer blood to points remote from the detector. Even efforts to avoid this situation can be defeated by movement of the patient, deliberate or inadvertent (e.g., the unconscious movement of a sleeping patient).\nStill another device for detecting leaks is described in U.S. Pat. No. 6,044,691. According to the description, the circuit is checked for leaks prior to the treatment operation. For example, a heated fluid may be run through the circuit and its leakage detected by means of a thermistor. The weakness of this approach is immediately apparent: there is no assurance that the system's integrity will persist, throughout the treatment cycle, as confirmed by the pre-treatment test. Thus, this method also fails to address the entire risk.\nYet another device for checking for leaks in return lines is described in U.S. Pat. No. 6,090,048. In the disclosed system, a pressure signal is sensed at the access and used to infer its integrity. The pressure wave may be the patient's pulse or it may be artificially generated by the pump. This approach cannot detect small leaks and is not very sensitive unless powerful pressure waves are used; in which case the effect can produce considerable discomfort in the patient.\nClearly detection of leaks by prior art methods fails to reduce the risk of dangerous blood loss to an acceptable level. In general, the risk of leakage-related deaths increases with the decrease in medical staff per patient driven by the high cost of trained staff. Currently, with lower staffing levels comes the increased risk of unattended leaks. Thus, there has been, and continues to be, a need in the prior art for a foolproof approach to detection of a return line leak or disconnection.\nIn an area unrelated to leak detection, U.S. Pat. No. 6,177,049 B1 suggests the idea of reversing the direction of blood flow for purposes of patency testing and access-clearing. The patency tests alluded to by the '049 patent refer simply to the conventional idea of forcing blood through each access to clear occlusions and to ascertain the flow inside a fistula."}
{"text": "This section introduces aspects that may be helpful in facilitating a better understanding of the inventions. Accordingly, the statements of this section are to be read in this light and are not to be understood as admissions about what is in the prior art or what is not in the prior art.\nFourth generation (4G) wireless mobile telecommunications technology, also known as Long Term Evolution (LTE) technology, was designed to provide high capacity mobile multimedia with high data rates particularly for human interaction. Next generation or fifth generation (5G) technology is intended to be used not only for human interaction, but also for machine type communications in so-called Internet of Things (IoT) networks.\nWhile 5G networks are intended to enable massive IoT services (e.g., very large numbers of limited capacity devices) and mission-critical IoT services (e.g., requiring high reliability), improvements over legacy mobile communication services are supported in the form of enhanced mobile broadband (eMBB) services providing improved wireless Internet access for mobile devices.\nIn an example communication system, user equipment (5G UE in a 5G network or, more broadly, a UE) such as a mobile terminal (subscriber) communicates over an air interface with a base station or access point referred to as a gNB in a 5G network. The access point (e.g., gNB) is illustratively part of an access network of the communication system. For example, in a 5G network, the access network is referred to as a 5G System and is described in 5G Technical Specification (TS) 23.501, V15.2.0, entitled “Technical Specification Group Services and System Aspects; System Architecture for the 5G System,” the disclosure of which is incorporated by reference herein in its entirety. In general, the access point (e.g., gNB) provides access for the UE to a core network (CN), which then provides access for the UE to other UEs and/or a data network such as a packet data network (e.g., Internet).\nTS 23.501 goes on to define a 5G Service-Based Architecture (SBA) which models services as network functions (NFs) that communicate with each other using representational state transfer application programming interfaces (Restful APIs).\nFurthermore, 5G Technical Specification (TS) 33.501, V15.1.0, entitled “Technical Specification Group Services and System Aspects; Security Architecture and Procedures for the 5G System” and 5G Technical Report (TR) 33.899, V1.3.0, entitled Technical Specification Group Services and System Aspects; Study on the Security Aspects of the Next Generation System,” the disclosures of which are incorporated by reference herein in their entireties, further describe security management details associated with a 5G network.\nSecurity management is an important consideration in any communication system. For example, protections in a 5G network against false bases stations, e.g., communication equipment or devices, operated by malicious actors, pretending to be legitimate base stations in a serving network, are critical to prevent such malicious actors from acquiring sensitive subscriber information that would allow them to, inter alia, act like a legitimate UE to the 5G network."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical information processing apparatus which irradiates a focused light beam to an optical recording medium to record and/or reproduce information.\n2. Related Background Art\nVarious disk, card and tape media which record information by light and reproduce the recorded information have been known.\nFor example, in an optical information processing apparatus which uses an optical disk, the optical disk is scanned by a light beam which is modulated by recording information and focused to a fine spot, and information is recorded as optically detectable record pit tracks (information tracks). In order to exactly record the information without difficulty, such as crossing of the information tracks, it is necessary to control the irradiation position of the light spot on a plane of the optical disk in a direction perpendicular to the scan direction (auto-tracking, hereinafter referred to as AT). Further, in order to irradiate the light spot as a fine spot which is stable in spite of warp or mechanical tolerance of the optical disk, it is necessary to control the irradiation position normal to the plane of the optical disk (auto-focusing, hereinafter referred to as AF). In a reproduction mode, both AT and AF are necessary, too. Various techniques for AT and AF have been known. Usually, a focusing error signal and a tracking error signal are derived from differences between photo-sensing planes of focusing and tracking photo-detectors, and an objective lens is driven by AT and AF actuators.\nIn the above apparatus, when recording or reproduction is to be started, the focusing control means is switched from an inactive state to an active state, that is, so-called focus pull-in is effected. In the focus pull-in, a focus control loop is opened and a triangular wave signal is applied to a focus actuator to move the objective lens up and down, as described in U.S. Pat. No. 4,542,491. When the objective lens reaches a position corresponding to a linear region of a focus error signal, the loop is closed. During the focus pull-in, the tracking control loop is kept open.\nIn the above method, since no drive force along the tracking path is applied to the objective lens in the focus pull-in, the objective lens may deviate from an intended mechanical center (a center of a movable range by the tracking actuator, at which an optical axis of the light beam and an optical axis of the objective lens coincide). Such a deviation may be caused by a residual strain of a spring which movably supports the objective lens or by an offset of an input circuit to the tracking actuator. If such a deviation occurs, the operation range may be unbalanced when the tracking actuator is activated after the focus pull-in. For example, assuming that the optical system of the apparatus assures an optical output in a range of .+-.250 .mu.m from the optical center and that there occurs a deviation of +100 .mu.m, the assurance of the optical output is in a range of -350 .mu.m to +100 .mu.m. Accordingly, if the light spot is track-jumped in a positive direction by 100 tracks (1.6 .mu.m/track) after the focus pull-in, the light beam is out of the range of assurance of the optical output."}
{"text": "Most conventional motor vehicles, such as the modern-day automobile, include a powertrain (sometimes referred to as “drivetrain”) that is generally comprised of an engine that delivers driving power through a multi-speed power transmission to a final drive system, such as a rear differential, axle, and wheels. Automobiles have traditionally been powered solely by a reciprocating-piston type internal combustion engine (ICE) because of its ready availability and relative cost, weight, and efficiency. Such engines include 4-stroke compression-ignited diesel engines and 4-stroke spark-ignited gasoline engines.\nHybrid vehicles, on the other hand, utilize alternative power sources to propel the vehicle, minimizing reliance on the engine for power, thereby increasing overall fuel economy. A hybrid electric vehicle (HEV), for example, incorporates both electric energy and chemical energy, and converts the same into mechanical power to propel the vehicle and power the vehicle systems. The HEV generally employs one or more electric machines that operate individually or in concert with an internal combustion engine to propel the vehicle. Since hybrid vehicles can derive their power from sources other than the engine, engines in hybrid vehicles can be turned off while the vehicle is propelled by the alternative power source(s).\nSeries hybrid architectures, sometimes referred to as Range-Extended Electric Vehicles (REEVs), are generally characterized by an internal combustion engine in driving communication with an electric generator. The electric generator, in turn, provides power to one or more electric motors that operate to rotate the final drive members. In effect, there is no direct mechanical connection between the engine and the drive members in a series hybrid powertrain. The lack of a mechanical link between the engine and wheels allows the engine to be run at a constant and efficient rate, even as vehicle speed changes—closer to the theoretical limit of 37%, rather than the normal average of 20%. The electric generator may also operate in a motoring mode to provide a starting function to the internal combustion engine. This system may also allow the electric motor(s) to recover energy from slowing the vehicle and storing it in the battery by regenerative braking.\nParallel hybrid architectures are generally characterized by an internal combustion engine and one or more electric motor/generator assemblies, each of which has a direct mechanical coupling to the power transmission. Most parallel hybrid designs combine a large electrical generator and a motor into one unit, providing tractive power and replacing both the conventional starter motor and the alternator. One such parallel hybrid powertrain architecture comprises a two-mode, compound-split, electro-mechanical transmission which utilizes an input member for receiving power from the ICE, and an output member for delivering power from the transmission to the driveshaft. First and second motor/generators operate to rotate the transmission output shaft. The motor/generators are electrically connected to an energy storage device for interchanging electrical power between the storage device and the first and second motor/generators. A control unit is provided for regulating the electrical power interchange between the energy storage device and motor/generators, as well as the electrical power interchange between the first and second motor/generators.\nRegardless of architecture, most hybrid powertrains generate driveline vibrations during normal operation, which range from imperceptible to unpleasantly noticeable. Significant driveline vibrations may be objectionable to a vehicle operator, and may reduce service life of the driveline components. Historically, driveline vibrations are mitigated by implementing systems which operate to cancel torque oscillations at one specific frequency, over a range of frequencies, or a set of frequencies chosen based upon gear ratio at which the driveline is currently operating. Such torque cancellation systems typically pass driveline inputs through signal conditioning filters, which may slow system responsiveness. Slow system response often leads to a “bump” or “overshoot” that occurs when there is an aggressive operator torque request, due to delays in transient responses required to develop filters.\nSome systems use a single feedback variable, typically engine speed, and command a single control signal, typically engine torque. However, single feedback/single control vibration control systems do not provide adequate damping in a system having multiple devices operable to generate vibrations in the driveline. As such, other systems employ a multivariate feedback control approach to provide active driveline damping for a hybrid powertrain. This approach provides dynamic coordination of all torque commands to control the transient response of the driveline using the hybrid transmission, including engine torque commands, electric motor torque commands, and clutch torque commands, as well as other controllable torque inputs."}
{"text": "1. Field of the Invention\nThis invention relates to stepping motors, and more particularly to a method and apparatus for determining the speed and position values of such motors.\n2. Description of the Prior Art\nThe stepping motor is increasingly used for applications calling for the motor to be operated quasi-continuously, because of its relatively simple digital control. The open loop mode and the feed-back mode are two operating modes for applications where the stepping motor is operated quasi-continuously. The open loop mode is highly susceptible to oscillations and the feed-back mode necessitates a feed-back signal indicating the respective operational stage of the motor shaft. Optically coded disc, attached to the motor shaft, provide this feed-back signal, by emitting one pulse for each motor step. To ensure that this pulse is accurate within a few percent of the step width, the discs pattern and the adjustment of the axle must meet very stringent requirements. An increase in the degree of accuracy of the disc is possible by using larger discs, but that would mean an increase in the rotational mass of the motor. The protection needed against contamination, the light source, and the light detector required render this method very expensive.\nWith each pulse a coded disc of the kind described provides information on the position of the motor shaft. Information on the speed can be obtained only as a function of the time difference between two feed-back pulses. Whenever this method is employed the speed during a step just completed can only be determined afterwards. This is a serious disadvantage when the motor moves very slowly shortly before it reaches its stand-still position. In other words, instantaneous speed data is not available.\nWhere a stepping motor is used to control the operational sequence the greatest problem is to accurately control the motor into its end position."}
{"text": "1. Field of the Invention\nThe present invention relates to improvement in or relating to a photoelectric conversion device in which a number of semiconductor elements are sequentially arranged on a substrate in side-by-side relation and connected in series. The invention also pertains to a method for the manufacture of such a photoelectric conversion device.\n2. Description of the Prior Art\nThere has been proposed in U.S. Pat. No. 4,315,096 a photoelectric conversion device of the type wherein a plurality n (n being an integer greater than one) of semiconductor elements U.sub.i to U.sub.n are sequentially formed side by side on a substrate having an insulating surface and are connected in series one after another.\nAccording to this semiconductor photoelectric conversion device, the semiconductor element U.sub.i (i=1, 2, . . . n) has a first electrode E.sub.i formed on the substrate, a non-single-crystal semiconductor laminate member Q.sub.i formed on the first electrode E.sub.i to form at least one semiconductor junction and a second electrode F.sub.i formed on the non-single-crystal semiconductor laminate member Q.sub.i in opposing relation to the first electrode E.sub.i. The second electrode F.sub.j+1 of the semiconductor element U.sub.j+1 (J=1, 2, . . . (n-1)) is coupled with the first electrode E.sub.j of the semiconductor element U.sub.j through an extension K.sub.j of the second electrode F.sub.j+1.\nIn such a photoelectric conversion device, in order to prevent lowering of its photoelectric conversion efficiency, it is necessary that the non-single-crystal semiconductor laminate member Q.sub.i and the second electrode F.sub.i be held in good contact with each other for a long period of time.\nIn the photoelectric conversion device of the abovesaid U.S. patent, however, no particular attention is paid to such a structure that ensures retention of good contact between the non-single-crystal semiconductor laminate member Q.sub.i and the second electrode F.sub.i.\nAccordingly, this conventional photoelectric conversion device has the defect that high photoelectric conversion efficiency cannot be maintained for a long period of time.\nFurther, it is described in the abovesaid U.S. patent that the second electrode F.sub.i is formed by a conductive layer through laser beam scanning. But it is not taken into account that during the laser beam scanning a conductive material forming the conductive layer enters into the non-single-crystal semiconductor laminate member to impair the electrical insulation between the second electrodes F.sub.j and F.sub.j+1. Accordingly, the electrical insulation between the second electrodes F.sub.j and F.sub.j+1 is poor.\nTherefore the photoelectric conversion device of the abovesaid U.S. patent has the defect of low photoelectric conversion efficiency."}
{"text": "The present invention relates to semiconductor integrated circuits containing memory arrays, and particularly those arrays incorporating array lines having extremely small pitch, and more particularly those having a three-dimensional memory array.\nSemiconductor integrated circuits have progressively reduced their feature linewidths into the deep sub-micron regime. Moreover, recent developments in certain memory cell technologies have resulted in word lines and bit line having an extremely small pitch. For example, certain passive element memory cell arrays may be fabricated having word lines approaching the minimum feature size (F) and minimum feature spacing for the particular word line interconnect layer, and also having bit lines approaching the minimum feature width and minimum feature spacing for the particular bit line interconnect layer. Moreover, three-dimensional memory arrays having more than one plane of memory cells have been fabricated containing such so-called 4F2 memory cells on each memory plane. Exemplary three-dimensional memory arrays are described in U.S. Pat. No. 6,034,882 to Johnson, entitled “Vertically Stacked Field Programmable Nonvolatile Memory and Method of Fabrication.”\nHowever, the area required for implementing decoder circuits for word lines and bit lines has not achieved such dramatic reductions. Consequently, interfacing the word line decoders and bit line decoders to such tightly spaced word lines and bit lines within such very dense arrays has become extremely difficult, and limits the density of memory arrays otherwise achievable. There remains a continued need for improved decoder structures capable of interfacing with large numbers of array lines having a very small pitch, and particularly if such array lines exist on more than one layer, as in a three-dimensional memory array having more than one plane or level of memory cells.\nAdditionally, integrated circuits incorporating a passive element memory array require a high-voltage and high-current programming voltage source due to the large number of leakage paths in the array and the high voltage required to program the element conductivity. The leakage current represents a significant portion of the power dissipation of such circuits during programming. There remains a need for improved performance of such circuits, reduced leakage currents when writing, and faster write time of a selected memory cell."}
{"text": "Refrigeration installations comprising a closed circuit in which a heat-transfer fluid (for example glycol water) is forced to circulate between capsules filled with phase-change material and stacked in a tank (made of steel or of concrete) and then led to the zone that is to be cooled (the technology referred to as “encapsulated PCM”) are known. The thinner the wall of the capsules the better the coefficient of heat transfer from the capsules to the heat-transfer fluid. During the phase referred to as the “store-charging” phase, the heat-transfer fluid, cooled by a refrigeration compressor, circulates through the tank at a temperature lower than the temperature at which the phase-change material contained in the capsules changes state and this has the effect of solidifying the phase-change material contained in the capsules and therefore of storing a certain amount of refrigeration energy. During the phase referred to as the “store discharging” phase, the heat-transfer fluid circulates through the tank and, upon contact with the capsules filled with the solidified PCM, picks up stored refrigeration energy and transfers it to the zone that is to be cooled. This circulation causes the phase-change material in the capsules to melt progressively, which means that the phase-change material has to be returned periodically to the solid state (something which is done during the store-charging phases).\nCertain phase-change materials, notably water, occupy a greater volume in the solid state than in the liquid state and it is important for the capsule to be able to absorb this increase in volume without sustaining damage. One immediate solution is to partially fill the capsule with the phase-change material, the remainder of the volume being occupied by air and forming a free volume that can be gradually occupied as the phase-change material solidifies, at the expense of an increase in pressure in the capsule. Although simple to implement, this solution has the disadvantage of causing the thin wall of the capsule to stretch in a zone of weaker strength, and of allowing the casing of the capsule to deform by forming a dished shape in a zone of lesser strength of the casing under the effect of the pressure of the heat-transfer fluid.\nThe repeated nature of these deformations causes weakening of the casing which may ultimately yield.\nDocument FR2609536 describes a capsule completely filled with phase-change material and comprising a flexible casing which has hollow dished shapes that can be pushed back by the phase-change material as it solidifies, thereby allowing an increase in the internal volume of the capsule. As before, the repeated nature of the deformations of the casing ultimately weakens the latter.\nDocument FR2732453 itself describes a capsule with a thin and rigid casing containing a spherical absorption body held at the center of the capsule and occupying part of the internal volume of the capsule. The internal volume of the capsule, which is therefore decreased by the volume of the absorption body itself, is completely filled with phase-change material. The absorption body is compressible and therefore able to be compressed by the phase-change material as it solidifies so as to absorb the increase in volume of said material. The spherical shape of the expansion body associated with the spherical shape of the capsule leads the phase-change material to solidify from the periphery toward the center of the capsule without flowing, thereby avoiding the creation of detrimental internal stresses. However, fitting such an absorption body at the center of the capsule and keeping it there are tricky."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of medical equipment for respiratory therapy and more specifically to the user interface for a ventilator used for monitoring and controlling the breathing of a patient.\n2. Description of the Related Art\nModern patient ventilators are designed to ventilate a patient's lungs with breathing gas, and to thereby assist a patient when the patient's ability to breathe on his own is somehow impaired. As research has continued in the field of respiration therapy, a wide range of ventilation strategies have been developed. For example, pressure assisted ventilation is a strategy often available in patient ventilators and includes the supply of pressure assistance when the patient has already begun an inspiratory effort. With such a strategy, it is desirable to immediately increase the pressure after a breath is initiated in order to reach a target airway pressure for the pressure assistance. This rise in pressure in the patient airway which supplies breathing gas to the patient's lungs allows the lungs to be filled with less work of breathing by the patent. Conventional pressure assisted ventilator systems typically implement a gas flow control strategy of stabilizing pressure support after a target pressure is reached to limit patient airway pressure. Such a strategy also can include programmed reductions in the patient airway pressure after set periods of the respiratory cycle in order to prepare for initiation of the next patient breath.\nAs patient ventilator systems and their various components, including sensors and control systems, have become more sophisticated, and more understanding is gained about the physiology of breathing and the infirmities and damage which form the requirements for respiratory therapy, the number of variables to be controlled and the timing and interrelationships between the parameters have begun to confront the caregiver with a daunting number of alternative therapeutic alternatives and ventilator settings. Also, in such a complex environment, the interface between the ventilator and the caregiver has often not been adaptable to the capabilities of the operator, thus running the chance of either limiting the choices available to a sophisticated user or allowing a relatively less sophisticated user to chose poorly from the alternatives presented. Thus, it would be beneficial if a ventilator interface guided the user through the setup or therapy modification process, illustrating the relationship between changes, preventing incorrect or dangerous settings and sounding alarms or other audible indications of invalid settings when something is about to done that exceeds limits, but also allowing the advanced and sophisticated user to gain access to the full range of ventilator capabilities through an interface which both presents the various parameters and allows the visualization of their relationships.\nClinical treatment of a ventilated patient often requires that the breathing characteristics of the patient be monitored to detect changes in the breathing patterns of the patient. Many modern ventilators allow the visualization of patient breathing patterns and ventilator function and the caregiver adjusts the settings of the ventilator to fine tune the respiratory strategy being performed to assist the patient's breathing. However, these systems have been, up until now, relatively difficult to use by the unsophisticated user unless a limited number of options are selected. For example, in one prior art system, only a single respiratory parameter may be altered at a time. Moreover, the various respiratory parameters must often be entered into the ventilator controller in a prescribed order, or, where no order is prescribed, certain orders of entry should be avoided, otherwise the intermediate state of the machine before entry of the remaining parameters may not be appropriate for the patient. This inflexible approach to ventilator setup requires additional time and training if the user is to quickly and efficiently use the ventilator in a critical care environment.\nPrevious systems have also been deficient in that it is often difficult to determine the underlying fault that has caused an alarms to be sounded, and what controls or settings should be adjusted to cure the problem causing the alarm. For example, prior alarm systems have consisted of nothing more than a blinking display or light with an alarm to alert the user that a problem existed. Similarly, many prior art systems provided only limited assistance to a user or technician in setting the parameters to be used during treatment. For example, if a technician attempted to enter a setting that was inappropriate for the patient because of body size or for some other reason, the only alarm provided may have been an auditory indication that the value was not permitted, but no useful information was provided to assist the technician in entering an appropriate setting.\nOne problem consistently presented by prior art ventilator control systems has been that the user interface has offered relatively little to guide and inform the user during the setup and use of the ventilator. Prior systems typically utilized a single visual display of the operating parameters of the ventilator and sensed patient parameters. Alternatively, prior systems may have numerous fixed numeric displays, certain of which may not be applicable during all ventilation therapies. Even when more than one display has been provided, users typically received limited feedback, if any, from the control system indicating the effect that changing one particular setting had on the overall respiratory strategy. If a parameter was to be adjusted, the display would change to display that particular parameter upon actuation of the appropriate controls, and allow entry of a value for that parameter. However, the user was provided with no visual cue as to how the change in the parameter value would effect the overall ventilation strategy, and thus had no assistance in determining whether the value entered for the parameter was appropriate for the patient.\nWhat has been needed and heretofore unavailable in patient ventilators is a user friendly graphic interface that provides for simultaneous monitoring and adjustment of the various parameters comprising a respiratory strategy. Such an interface would also preferably guide sophisticated users in implementing ventilation therapies, provide guidance on the relationships between parameters as they are adjusted, allow rapid return to safe operation in the event that an undesirable strategy was inadvertently entered, provide alarms that are easily understood and corrected and present all of the relevant information in an easily understood and graphic interface. The present invention fulfills these and other needs."}
{"text": "Aspects of the disclosure relate to computer hardware and software. In particular, one or more aspects of the disclosure generally relate to computer hardware and software that can be used to provide an Internet edge cleansing farm.\nCyber attacks have become more prevalent and have had increasingly negative consequences over today's modern Internet network. Organizations that service customer requests over the Internet at a large scale are candidates for cyber attackers of the modern age. In particular, denial of service attacks are readily used and cyber attackers have found some success with this strategy. Mitigation tactics have been used to deter cyber attackers from implementing denial of service attacks. However, at times, confidential information is passed between organizations and customers. Accordingly, there is a need to mitigate against cyber attacks when confidential information may be involved in the communications between organizations and customers. In addition, cyber attackers are highly agile and adapt well. Accordingly, there is also a need for an adaptable mitigation strategy against cyber attacks."}
{"text": "The present invention relates to a nematic liquid crystal composition and, more specifically, relates to a liquid crystal composition for active matrix (AM) mode and to a liquid crystal display device (LCD) using the liquid crystal composition.\nAn active matrix mode-liquid crystal display devices (AM-LCD) is attracting a great deal of attention as the likely winner among LCD\"\"s as it enable to display with extreme precision, and is applied for displays such as monitors, note-type personal computers, digital still cameras and digital video cameras. Characteristics required for the liquid crystal composition of AM-LCD are shown in the items (1) to (5) below.\n(1) A liquid crystal composition should show a nematic phase in the range of temperature as wide as possible in order to widen the usable temperature range of the liquid crystal display devices (elevating the upper temperature limit of a nematic phase as high as possible, and lowering the lower temperature limit of a nematic phase as low as possible).\n(2) Viscosity of the liquid crystal composition should be as low as possible in order to accelerate the response speed of the liquid crystal display device.\n(3) Optical anisotropy (xcex94n) of the liquid crystal composition should be able to have a proper value depending on the cell thickness (d) in order to enhance the contrast of the liquid crystal display device.\n(4) Values of specific resistivity (specific resistance) on the liquid crystal composition should be increased and the voltage holding ratio of the cell containing the liquid crystal composition should be large, in order to enhance the contrast of the liquid crystal display device. The voltage holding ratio should be large especially at a high temperature region.\n(5) Threshold voltage of the liquid crystal composition should be lowered in order to downsize the battery which is the power source to drive liquid crystal display device.\nRecently it has been requested keenly to apply for an animation in the LCD and various LCD display modes have been studied. The addressing mode of the AM-LCD employs TN display mode in which the orientation of liquid crystal molecules between the upper and lower substrates is twisted by 900xc2x0. In the TN display mode, it is necessary to keep the product, xcex94nxc2x7d, of optical anisotropy (xcex94n) of a liquid crystal material to be filled in a cell and cell thickness (d xcexcm) at a certain value (for example xcex94nxc2x7d=0.5 etc.) in order to obtain an optimum contrast and to avoid coloration by the interference of the liquid crystal cell when no voltage is applied, as reported by G.Bauer (Cryst. Liq., 63, 45 (1981)). Therefore, if the liquid crystal material having large xcex94n is used, the value of d can be decreased. Response speed (xcfx84) is proportional to viscosity (xcex7) of the liquid crystal material and to d squared, as proposed by E. Jakeman et al. (Phys. Lett. , A, 39 (1972) 69). When the liquid crystal material having large xcex94n is used, the thickness of cell constituting liquid crystal display device can be decreased, and thus xcfx84 can be increased. As such, the liquid crystal composition having large xcex94n and low viscosity is very useful for the liquid crystal display devices.\nAlso the development of the display intended for outdoor use has been requested with the increase of a portable type display. For being bearable in the outdoor use, materials are requested to have a nematic phase over the wide range of temperature exceeding the temperature range-of environment in use. To widen a nematic phase range of the liquid crystal composition, it is necessary to use liquid crystal compounds having a high clearing point and good miscibility with other liquid crystal compounds. Generally, a compound having a high clearing point, namely having many six member-rings in the chemical structure may be used for elevating the clearing point. However, miscibility at a low temperature region may tend to be a problem in such case.\nIt is also requested the liquid crystal composition with high reliability such as a high voltage holding ratio (V.H.R.) or large specific resistivity in order to keep a high contrast especially in the AM-LCD.\nBased on the background described above, the specification of WO 96/11897 discloses a novel liquid crystal compound having large dielectric anisotropy (xcex94n) together with extremely low viscosity and disclosed a liquid crystal composition comprising the compound for low voltage driving in various modes such as a AM mode or a super twisted nematic mode (STN mode). JP 10-251186 A describes a compound similar to the compound of formula (1-3) in the present invention as the compound having large xcex94xcex5 and small temperature dependence.\nThe liquid crystal composition for the AM-LCD disclosed in WO 96/11897 described above (the composition does not contain compounds having cyano in the terminal, and compounds having cyano can not be used for the liquid crystal composition of the AM-LCD because of its low voltage holding ratio) has drawbacks that the V.H.R. is low and optical anisotropy is low as is shown in the Comparative Example of the present invention.\nJP 10-251186 A discloses a composition using the compound similar to that of formulae (1-1) to (1-3) of the present invention, however, the similar compound has large and positive xcex94xcex5. When such compound is used, viscosity of the composition is increased. Then, the compound has a drawback of low response speed in the LCD. The similar compound has a drawback that xcex94n is small and that a high voltage holding ratio required for the composition of TFT can not be attained by the combination of a cyano compound.\nAlthough liquid crystal compositions were studied, the liquid crystal composition for AM-LCD has been required especially to keep a high voltage holding ratio at a high temperature region in order to enhance the contrast of the liquid crystal display device, to have a wide range of a liquid crystal phase in order to have a wide addressing temperature range, and to have low viscosity in order to accelerate response speed, while maintaining the characteristics (1) to (5) described above. Enlarging xcex94n has also been requested by the requirement being smaller cell gap.\nThe liquid crystal composition has been always required to be improved though enthusiastic efforts are made for various purposes.\nThe object of the present invention is to provide a liquid crystal composition especially having a high voltage holding ratio at a high temperature region in order to enhance the contrast of a liquid crystal display device, having a wide range of a liquid crystal phase which means a high clearing point and good miscibility at low temperature, having low viscosity in order to accelerate response speed of the liquid crystal display device, and having properly large xcex94n in order to enhance the contrast in the liquid crystal display device, while satisfying general characteristics required for the liquid crystal composition of the AM-LCD.\nThe present inventors have made enthusiastic efforts to achieve the above object, have found the liquid crystal composition described below can achieve the purpose of the present invention, and have completed the present invention.\nThe liquid crystal composition of the present invention is described in the items 1 to 4 below.\n1. A liquid crystal composition characterized by comprising at least one compound selected from the group of compounds expressed by formulae (1-1), (1-2) and (1-3) as a first component and comprising at least one compound selected from the group of compounds expressed by formulae (2-1), (2-2), (2-3), (2-4), (2-5), and (2-6) as a second component. \nxe2x80x83wherein R1 to R9 each independently represents alkyl or alkoxy having 1 to 10 carbon, alkenyl or alkoxymethyl having 2 to 10 carbon; R10 to R12 each independently represents alkyl or alkoxy having 1 to 10 carbon, alkenyl or alkoxymethyl having 2 to 10 carbon, F, Cl, CF3, OCF3, or OCF2H; X1 to X6 each independently represents F, CF3, OCF3, OCF2H or Cl; Y1 to Y11 each independently represents H or F; Z1 to Z4 each independently represents a single bond or xe2x80x94C2H4xe2x80x94.\n2. The liquid crystal composition described in the item 1 described above characterized by comprising 3 to 45% by weight of the first component, and comprising 25 to 97% by weight of a second component each based on the total weight of the liquid crystal composition.\n3. A liquid crystal composition comprising at least one compound selected from the group of compounds expressed by formulae (1-1) to (1-3) as a first component, comprising at least one compound selected from the group of compounds expressed by formulae (2-1) to (2-6) as a second component, and comprising at least one compound selected from the group of compounds expressed by formula (3) as a third component \nxe2x80x83wherein R13 and R14 each independently represents alkyl having 1 to 10 carbon.\n4. The liquid crystal composition in the item 3 described above characterized by comprising 3 to 45% by weight of the first component, comprising 25 to 97% by weight of the second component, and comprising 25% and less by weight of the third component each based on the total weight of the liquid crystal composition.\nThe liquid crystal display device of the present invention is described in the item 5 below.\n5. The liquid crystal display device composed by using the liquid crystal composition in any one of the items 1 to 4 described above.\nFollowings are explanation of preferable embodiments of compounds constituting the liquid crystal composition of the present invention.\nAmong the compounds expressed by formula (1-1) as the first component of the liquid crystal composition of the present invention, the compounds expressed by formulae (1-1-1) to (1-1-12) in the following are preferably used. In these formulae, R and Rxe2x80x2 each independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (1-2) as the first component of the liquid crystal composition of the present invention, the compounds expressed by formulae (1-2-1) to (1-2-12) in the following are preferably used. In these formulae, R and Rxe2x80x2 each independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (1-3) as the first component of the liquid crystal composition of the present invention, the compounds expressed by formulae (1-3-1) to (1-3-12) in the following are preferably used. In these formulae, R and Rxe2x80x2 each independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (2-1) as the second component of the liquid crystal composition of the present invention, the compounds expressed by formulae (2-1-1) to (2-1-30) in the following are preferably used. In these formulae, each R independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (2-2) as the second component of the liquid crystal composition of the present invention, the compounds expressed by formulae (2-2-1) to (2-2-40) in the following are preferably used. In these formulae, each R independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (2-3) as the second component of the liquid crystal composition of the present invention, the compounds expressed by formulae (2-3-1) to (2-3-20) are preferably used. In these formulae, each R independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (2-4) as the second component of the liquid crystal composition of the present invention, the compounds expressed by formulae (2-4-1) to (2-4-9) are preferably used. In these formulae, each R independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (2-5) as the second component of the liquid crystal composition of the present invention, the compounds expressed by formulae (2-5-1) to (2-5-9) in the following are preferably used. In these formulae, each R independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nAmong the compounds expressed by formula (2-6) as the second component of the liquid crystal composition of the present invention, the compounds expressed by formulae (2-6-1) to (2-6-9) are preferably used. In these formulae, each R independently represents alkyl or alkoxy having 1 to 10 carbon, or alkoxymethyl or alkenyl having 2 to 10 carbon. \nFollowings are explanation on the role of compounds constituting the liquid crystal composition of the present invention.\nThe compounds expressed by formulae (1-1) to (1-3), which are the first component of the liquid crystal composition of the present invention, has features of a very large xcex94n being about 0.13 to 0.25, of relatively small xcex7 being 20 to 60 mPaxc2x7s compared with other three or four ring-compounds, and of a large value of specific resistivity. For this reason, the compounds expressed by formulae (1-1) to (1-3) of the present invention are used for the purpose to enlarge xcex94n, to minimize xcex7, and to adjust threshold voltage, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region. As the compounds expressed by formulae (1-1) and (1-2) have four rings and have very high TNI (the upper temperature limit of a liquid crystal phase) being 100xc2x0 C. to 180xc2x0 C., the liquid crystal compositions having a high TNI can be prepared by the use of these four ring-compounds.\nThe compounds expressed by formulae (2-1) to (2-6), which are the second component of the present invention, have features of large xcex94n being about 0.03 to 0.18, of large xcex94xcex5 being about 6 to 32, and of a large value of specific resistivity. For this reason, the compounds expressed by formulae (2-1) to (2-6) of the present invention are used for the purpose to keep large xcex94n, to widen a range of a liquid crystal phase, and to adjust threshold voltage lower, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nFor further detailed explanation, the compounds expressed by formula (2-1), which are the second component of the present invention, have three rings and have xcex94n being about 0.03 to 0.10 and a large xcex94xcex5 value being about 6 to 13. They also have features of relatively small xcex7 and high specific resistivity. For this reason, the compounds expressed by formula (2-1) of the present invention are used for the purpose to decrease xcex7, and especially to adjust threshold voltage and xcex94n, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (2-2), which are the second component of the present invention, have three rings and have relatively large xcex94n being about 0.10 to 0.14 and a large value of xcex94xcex5 being about 9 to 18. They also have features of relatively small xcex7 and high specific resistivity. For this reason, the compounds expressed by formula (2-2) of the present invention are used for the purpose to minimize xcex7, to keep An relatively large, and especially to adjust threshold voltage being small, while maintaining a high voltage holding ratio of liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (2-3), which are the second component of the present invention, have four rings and have extremely a high TNI (the upper temperature limit of a liquid crystal phase) of 180xc2x0 C. or higher. They also have relatively large xcex94n being about 0.13 to 0.16 and relatively large xcex94xcex5 being about 10 to 14. Further they have features of relatively small xcex7 and high specific resistivity. For this reason, the compounds expressed by formula (2-3) of the present invention are used for the purpose to elevate the upper temperature limit of a liquid crystal phase of the liquid crystal composition, to keep xcex94n relatively large, and especially to adjust threshold voltage, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (2-4), which are the second component of the present invention, have three rings and have considerably large xcex94xcex5 being 15 to 32 and xcex94n being about 0.05 to 0.09. They also have features of a large value of specific resistivity. For this reason, the compounds expressed by formula (2-4) of the present invention are used especially to reduce threshold voltage, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (2-5), which are the second component of the present invention, have three rings and have considerably large xcex94xcex5 of about 15 to 32 and relatively a large value of xcex94n being about 0.1 to 0.15. They also have features of high specific resistivity. For this reason, the compounds expressed by formula (2-5) of the present invention are used especially to enlarge xcex94n and to reduce threshold voltage, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (2-6), which are the second component of the present invention, have four rings and have extremely high TNI being about 100xc2x0 C. to 130xc2x0 C., considerably large xcex94xcex5 being about 15 to 30 and relatively a large value of xcex94n being about 0.14 to 0.18. They also they have features of a large value of specific resistivity. For this reason, the compounds expressed by formula (2-6) of the present invention are used to elevate the upper temperature limit of a liquid crystal phase, and especially to enlarge xcex94n and to reduce threshold voltage, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (3), which are the third component of the present invention, have features of xcex94xcex5 being close to 0, of a large value of specific resistivity, of extremely large xcex94n of 0.2 or more which is extremely large, and of a high upper temperature limit of a nematic phase being 250xc2x0 C. or more. For this reason, the compounds expressed by formula (3) of the present invention are used for the purpose to elevate the upper temperature limit of a nematic phase, to adjust threshold voltage, and to enlarge xcex94n considerably, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nSummarizing the above, the liquid crystal composition of the present invention has features widening the range of a liquid crystal phase, especially minimizing xcex7 and enlarging xcex94n, while keeping a high upper temperature limit of the liquid crystal phase of the liquid crystal composition based on the first component, and further adjusting the range of the liquid crystal phase, xcex94n, and threshold voltage based on the second component.\nFollowings are explanation on preferable ratios in the contents of components which constitute the liquid crystal composition of the present invention.\nThe content of the first component in the liquid crystal composition of the present invention is preferably 3 to 45% by weight and the content of the second component therein is preferably 25 to 97% by weight each based on the total weight of the liquid crystal composition.\nFor further detailed explanation, the compounds expressed by formulae (1-1) to (1-3) are desirable to be mixed into the composition as much as possible to enlarge xcex94n and to minimize xcex7. However, they may elevate the lower temperature limit of a nematic phase of the liquid crystal composition if they are contained too much in the composition. For this reason, the compounds expressed by formulae (1-1) to (1-3), which are the first component of the present invention, are preferably 45% by weight and less based on the total weight of the liquid crystal composition. The ratio of the first component of the present invention in the liquid crystal composition is preferably 3% or more by weight to enlarge xcex94n, to widen the range of a liquid crystal phase, and to minimize xcex7, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formulae (2-1) to (2-6), which are the second component of the present invention, may lower the upper temperature limit of a nematic phase of the liquid crystal composition, increase xcex7, and lower xcex94n, when they are contained in the composition in large quantities. For this reason, the ratio of the compounds expressed by formulae (2-1) to (2-6), which are the second component of the present invention, is preferably 97% and less by weight based on the total weight of the liquid crystal composition. The ratio of the second component in the liquid crystal composition is preferably 25% or more by weight to keep low threshold voltage and to lower the lower temperature limit of a nematic phase, while maintaining a high voltage holding ratio of the liquid crystal composition at a high temperature region.\nThe compounds expressed by formula (3), which are the third component of the present invention, may elevate the lower temperature limit of a nematic phase of the liquid crystal composition and may increase the threshold voltage because their xcex94xcex5 are close to 0, when they are contained in the composition in large quantities. For this reason, the ratio of compounds expressed by formula (3) of the present invention is preferably 25% and less by weight based on the total weight of the liquid crystal composition.\nThe method of the preparation on the compounds expressed by formula (1-3) of the present invention, for example, that on the compounds expressed by formula (1-3-2) is described in JP 10-251186 A. On the compounds expressed by formulae (2-1), (2-2) and (2-3), for example, those expressed by formulae (2-1-16), (2-2-11) and (2-3-12), the method of their preparation is described in JP 2-233626 A. On the compounds expressed by formulae (2-4), (2-5) and (2-6) of the present invention, for example, those expressed by formulae (2-4-5), (2-5-5) and (2-6-5), the method of their preparation is described JP 10-251186 A. On the compounds expressed by formula (3), the method of their preparation is described in JP 2-237949 A. Then, each compound constituting the composition of the present invention can be prepared based on prior arts.\nTo the liquid crystal composition of the present invention, liquid crystal compounds other than those expressed by formulae described above can be mixed in the amount not affecting the object of the present invention. The liquid crystal composition used in the present invention can be prepared by the conventional methods. In general, the method is that various compounds are mixed and dissolved each other at a high temperature. To the liquid crystal composition of the present invention, chiral doping agents such as cholesteryl nonanoate (CN) or CM-43L expressed by the following formula may be added in order to adjust to the required twist angle by inducing the spiral structure to liquid crystal molecules. \nThe liquid crystal composition of the present invention can also be used for a guest-host mode by adding dichroic dyes such as phthalocyanine type, styril type, azo type, azomethine type, azoxy type, quinophthalone type, anthraquinone type and tetrazine type. It can also be used for a polymer-dispersed type liquid crystal display device, a birefringence control mode or a dynamic scattering mode. It can also be used for in-plane switching mode. Further, it is preferable for OCB mode utilizing its large xcex94n."}
{"text": "Thermal-magnetic trip units used within residential and commercial molded case circuit breakers are generally limited by geometric considerations from providing low current magnetic trip response. U.S. Pat. No. 4,513,268 describes a residential type molded case circuit breaker incorporating a thermal-magnetic trip unit in accordance with the prior art. U.S. Pat. No. 4,951,015 describes a movable core that is designed to move into the gap existing between the core and armature of a magnetic trip unit to reduce the primary air gap and increase the magnetic flux. The movable core effectively allows the circuit breaker to trip at lower current levels. U.S. Pat. Nos. 3,179,767, 3,278,707 and 3,278,708 each describe the use of an additional turn of wire around the magnet used within the thermal-magnetic trip unit to increase the magnetic forces on the armature at low currents.\nAdditionally, U.S. patent application Ser. No. 841,182 entitled \"Thermal-Magnetic Trip Unit with Low Current Response\" now U.S, Pat. No. 5,173,674 describes a molded case circuit breaker trip unit employing a pivotally-arranged magnet that moves unit employing a pivotally-arranged magnet that moves toward the armature to reduce the magnetic gap separation distance.\nThe aforementioned thermal-magnetic trip assemblies are found to add to the materials and assembly costs of the residential circuit breakers employing thermal-magnetic trip units.\nThe addition of supplemental magnets and armatures correspondingly increases the manufacturing tolerances that must be carefully controlled to insure compliance with the relevant industry standards.\nOne purpose of the invention accordingly is to provide a thermal-magnetic trip unit providing low current magnetic trip response with automatic tolerance compensation at relatively low cost."}
{"text": "1. Field of the Invention\nThis invention generally relates to anchoring arrangements for floating structures and more particularly to anchoring arrangements for floating structures such as converted tankers.\n2. Description of the Prior Art\nThe traditional anchoring arrangements comprising any of the known types of anchors cater to the anchoring requirements of seagoing vessels and other floating structures, for example, tankers, which are either at a dock or a harbor. Occasionally, certain conventional anchoring arrangements are used to anchor vessels in the open sea. However, it is known that conventional anchoring will not be very effective if the floating structure or vessel is in high seas, especially when it is required to connect a transfer hose or an underwater transfer tube carrying a fluid, for example oil or gas, to the floating vessel. Even with anchoring, a certain degree of shifting of the floating structure, or even a certain degree of rotation of the floating structure, might occur in high seas. At least one consequence of such shifting or rotation is that the transfer hose, for example, tends to get damaged or even disconnected, causing serious consequences. The problem is especially serious in the case of ocean vessels such as converted tankers in high seas wherein there is need for having at least a transfer tube constantly connected to the tanker through the ocean water, invariably to the sea bed or the shore or another vessel. Certain structural arrangements have been used heretofore to cater to the needs of stably and permanently anchoring converted tankers, simultaneously making provision for a transfer pipe or tube connected between the tanker and the sea bed. However, prior art arrangements have always required a significant mass of steel to fabricate the underwater anchoring structure, consequently rendering the equipment expensive from the poin the view of installation and maintenance. Furthermore, it has been found that it is desirable to provide an underwater means which permits relative rotary movement between the tanker superstructure and the anchoring arrangement to make a provision for oscillatory and rotary swaying movements of the tanker superstructure in high seas. Preferably, the underwater means should have a pivot bearing mechanism to easily permit such movement. Also, it is desirable that there should be easy access for maintenance personnel to reach at least part of the bearing mechanism to attend to maintenance and replacement work, preferably in a dry atmosphere. Such facilities are not available in any known prior art anchoring arrangement.\nThere is, therefore, a great need for an anchoring arrangement devoid of the disadvantages and limitations of prior art and including the more desirable features which are discussed above.\nThere is also a great need for an anchoring and transfer system which is low in cost and can be installed in a converted tanker using only minimum labor, wherein only a small working area is required at the anchorage. There is also a need for an anchoring arrangement including a transfer system wherein there is excellent accessibility to the anchoring apparatus, and wherein the fluid transfer system is weatherproof, and protected against collision."}
{"text": "Minerals are essential elements for the growth of all organisms. Dietary minerals can be derived from many source materials, including plants. For example, plant seeds are a rich source of minerals since they contain ions that are complexed with the phosphate groups of phytic acid molecules. These phytate-associated minerals may, in some cases, meet the dietary needs of some species of farmed organisms, such as multi-stomached ruminants. Accordingly, in some cases ruminants require less dietary supplementation with inorganic phosphate and minerals because microorganisms in the rumen produce enzymes that catalyze conversion of phytate (myo-inositol-hexaphosphate) to inositol and inorganic phosphate. In the process, minerals that have been complexed with phytate are released. The majority of species of farmed organisms, however, are unable to efficiently utilize phytate-associated minerals. Thus, for example, in the livestock production of monogastric animals (e.g., pigs, birds, and fish), feed is commonly supplemented with minerals and/or with antibiotic substances that alter the digestive flora environment of the consuming organism to enhance growth rates.\nAs such, there are many problematic burdens—related to nutrition, ex vivo processing steps, health and medicine, environmental conservation, and resource management—that are associated with an insufficient hydrolysis of phytate in many applications. The following are non-limiting examples of these problems:                1) The supplementation of diets with inorganic minerals is a costly expense.        2) The presence of unhydrolyzed phytate is undesirable and problematic in many ex vivo applications (e.g. by causing the presence of unwanted sludge).        3) The supplementation of diets with antibiotics poses a medical threat to humans and animals alike by increasing the abundance of antibiotic-tolerant pathogens.        4) The discharge of unabsorbed fecal minerals into the environment disrupts and damages the ecosystems of surrounding soils, fish farm waters, and surface waters at large.        5) The valuable nutritional offerings of many potential foodstuffs remain significantly untapped and squandered.        \nConsequently, phytate-containing foodstuffs require supplementation with exogenous nutrients and/or with a source of phytase activity in order to amend their deficient nutritional offerings upon consumption by a very large number of species of organisms.\nConsequently, there is a need for means to achieve efficient and cost effective hydrolysis of phytate in various applications. Particularly, there is a need for means to optimize the hydrolysis of phytate in commercial applications. In a particular aspect, there is a need to optimize commercial treatment methods that improve the nutritional offerings of phytate-containing foodstuffs for consumption by humans and farmed animals."}
{"text": "Macronutrients (N, K, Ca, Mg, P, and S) and Micronutrients (Fe, B, Mn, Zn, Cu, Mo, Co, and Ni) are crucial to a plant's growth, development, disease resistance and various metabolic pathways such as photosynthesis. Plant available micronutrient insufficiencies are due to traditional farming methods that have exhausted the soil and to the micronutrient metals existing as water insoluble salts and complexes. Many of the water insoluble forms in the soil involve a metal cation and boron, sulfur, or phosphorous based anions. A deficiency in micronutrients results in poor plant growth and development and thus in diminished yields (Mortvedt 1990). Plant requirements for many of the micronutrients can be as low as parts/million in the plant tissue. It is known that increasing the plant available micronutrient metal ions by addition of complexed metal ions to the soil or to plant foliage or by freeing up micronutrients, bound in the soil as an insoluble salts or complexes, in a plant absorbable form can help to significantly alleviate soil deficiencies and assist in development, growth, and disease resistance of the plants.\nPhosphorous is second to nitrogen as the most limiting macronutrient. In the case of phosphorus fertilizer, 40% of landscape soil is considered to contain inadequate levels of phosphorus for woody plant growth. Moreover, most of the phosphorus in the soil is largely inaccessible as it exists in a form that is not soluble in water and thus is not readily available to plants. In some cases, only 0.01% of the total soil phosphorus is in the form of a water soluble ion, the only form which can be absorbed by the plant. Adequate and accessible soil phosphorus is essential for optimal crop yields. Phosphorus enables a plant to store and transfer energy, promotes root, flower and fruit development, and allows early maturity. Phosphorus is also involved in many processes critical to plant development such as photosynthesis where plants utilize organic phosphorous compounds when converting sunlight to energy. Without enough phosphorus present in the soil, plants cannot grow sufficient root structure, which is key to the plant's ability to absorb water and nutrients from the soil. Moreover, woody plants, without sufficient root structure cannot maintain an equilibrium between roots and shoots, which is key to surviving drought, windy weather, and/or pests. Many of the nutrients required by plants are locked into salts and complexes that are water insoluble and therefore not plant available. To overcome these challenges, the agriculture industry has turned to chelates and anionic based polymers to form water soluble complexes with metal cations such as the micronutrients Ca, Mg, Mn, Fe, Cu, Co Ni, Zn, and Mo resulting in freeing up bound macronutrients such as phosphorous. The current delivery system technology of the chelates and polymer based products is water. Water is not only an excellent solubilizing/dispersing medium for chelates and [P(OA)]s, but can solvate a high load of water soluble metal salts. However, the use of water soluble metal salts can form insoluble complexes with chemistries that allow them to be available in the soil but unavailable to the plant.\nCoating a fertilizer with water based products can result in severe clumping of the fertilizer granules during blending, or gelling of the [P(OA)]s due to high electrolyte content caused by the fertilizer granule dissolving into the water. Clumping has a negative impact on its effectiveness to complex with metal cations, and/or it requires a drying step for seed coatings to prevent pre-mature sprouting or the growth of mold and mildew, which ultimately destroys the seeds. The use of aqueous based systems also has a deleterious impact on the urease inhibitor NBPT. The agricultural industry needs a technology that is able to easily, safely, evenly, and economically coat fertilizer granules and seeds with non-aqueous, liquid formulations that contain [P(OA)]s that can form water soluble metal cation complexes and free up bound macronutrients such as phosphorous."}
{"text": "A furnace, a reaction column, a heat exchanger, and so on of, for example, a hydrodesulfurization apparatus are exposed to fluids containing high-temperature sulfides during the operation, whereby iron sulfide is formed on the surface thereof. This iron sulfide, when exposed to the air, is hydrolyzed by the action of oxygen and moisture and is converted into polythionic acid, causing the occurrence of stress-corrosion cracking of an austenitic stainless steel used in the equipment.\nIn order to eliminate this problem, a method in which in stopping the operation, the fluids are withdrawn from the equipment and the inside of the equipment is washed and neutralized with an aqueous solution of an inorganic alkali such as sodium carbonate, caustic soda, or ammonia has heretofore been employed (see NACE Standard, RP01-70, titled \"Protection of Austenitic Stainless Steel in Refineries Against Stress Corrosion Cracking by Use of Neutralizing Solutions During Shut Down\").\nIn accordance with the above method comprising washing and neutralizing with an aqueous alkali solution, however, because the surface of the equipment is wet with fluids containing sulfides and repels the aqueous alkali solution, contact of the aqueous alkali solution with iron sulfide formed on the surface of the equipment is achieved insufficiently such that protection from the occurrence of stress-corrosion cracking cannot be ensured. Furthermore, the aqueous alkali solution for washing and neutralization sometimes remains in dead portions of the equipment and pipes to cause corrosion. Moreover, the above method involves such a problem that it is necessary to once withdraw the fluid remaining in the equipment and then introduce the aqueous alkali solution, which makes the operation complicated."}
{"text": "This invention relates to an exhaust gas recirculation (EGR) system for use in a diesel engine for recirculating a controlled rate of exhaust gases to the engine and, more particularly, to such an EGR system employing an electronic control unit adapted to calculate a target EGR ratio based upon engine operating parameters for controlling the rate of exhaust gases recirculated to the engine so as to obtain the calculated target EGR ratio.\nIn order to minimize emission of noxious pollutants discharged from a diesel engine to the atmosphere, it is the current practice to suppress combustion by recirculating a controlled rate of exhaust gases to the engine through an EGR passage having therein an EGR valve and connecting the engine exhaust passage to the engine intake passage downstream of the throttle valve. The rate of exhaust gases recirculated to the engine, which has a significant effect on both emission of nitrogen oxides and production of carbon fine particles, is determinative on not only the position of the EGR valve but also the position of the throttle valve across which a pressure differential exists in aid of introducing exhaust gases into the intake passage from the EGR passage. For example, the rate of exhaust gas flow through the EGR passage increases as the EGR valve moves in an opening direction for the same throttle valve position or as the throttle valve moves in a closing direction for the same EGR valve position. The position of the throttle valve, which also determines the rate of air flow to the engine, should be controlled properly to maintain optimum engine output performance in accordance with engine operating conditions.\nExhaust gas recirculation (EGR) control systems are well-known which involve an electronic control unit for providing accurate EGR ratio open-loop control. Such an electronic control unit calculates a target value for the EGR ratio meeting with requirements relating to engine output and exhaust performances as close as possible based upon engine operating parameters such as engine speed, accelerator pedal depression (indicated by fuel injection pump control sleeve or control rack position), fuel injection timing, engine coolant temperature, engine oil temperature, and the like and controls the EGR valve and throttle valve positions to obtain the calculated target EGR ratio. Such EGR ratio open-loop control has the distinct advantage in extremely fast response to engine operating condition changes.\nIn case where a deviation occurs between the calculated target EGR ratio value and the actual EGR ratio requirement due to errors in measurement in making and assembling engine parts such as the EGR valve and the throttle valve, changes in the engine part characteristics with the passage of time caused by mechanical wear and accumulated carbons on the engine parts, and the like, however, the EGR open-loop control system cannot correct the target EGR ratio value for the deviation. In addition, with an emission control device such as a soot collector located in the engine exhaust system for purifying engine exhaust emissions, the EGR open-loop control system cannot provide accurate EGR ratio control due to exhaust pressure changes caused by soot collected in the emission control device.\nThe present invention provides an improved and novel exhaust gas recirculation control system which open-loop controls the EGR ratio in accordance with a target EGR ratio value calculated based upon engine operating parameters and which, each time the vehicle travels a predetermined distance, calculates an actual EGR ratio value based upon measurements of the rate of air flow to the engine and the rate of exhaust gases recirculated to the engine and corrects the calculated target EGR ratio value for the deviation between the actual and target EGR ratio values, thereby eliminating the limitations and drawbacks inherent in previous EGR open-loop control systems."}
{"text": "Alignment characteristics of magnetic fields have been used to achieve precision movement and positioning of objects. A key principle of operation of an alternating-current (AC) motor is that a permanent magnet will rotate so as to maintain its alignment within an external rotating magnetic field. This effect is the basis for the early AC motors including the “Electro Magnetic Motor” for which Nikola Tesla received U.S. Pat. No. 381,968 on May 1, 1888. On Jan. 19, 1938, Marius Lavet received French Patent 823,395 for the stepper motor which he first used in quartz watches. Stepper motors divide a motor's full rotation into a discrete number of steps. By controlling the times during which electromagnets around the motor are activated and deactivated, a motor's position can be controlled precisely. Computer-controlled stepper motors are one of the most versatile forms of positioning systems. They are typically digitally controlled as part of an open loop system, and are simpler and more rugged than closed loop servo systems. They are used in industrial high speed pick and place equipment and multi-axis computer numerical control (CNC) machines. In the field of lasers and optics they are frequently used in precision positioning equipment such as linear actuators, linear stages, rotation stages, goniometers, and mirror mounts. They are used in packaging machinery, and positioning of valve pilot stages for fluid control systems. They are also used in many commercial products including floppy disk drives, flatbed scanners, printers, plotters and the like.\nMoreover, commercial, consumer, and industrial products and fabrication processes abound with a myriad of fasteners, latches, hinges, pivots, bearings and other devices that are conventionally based on mechanical strength and shape properties of materials rather than magnetic field properties because the magnetic field properties have been inadequate or otherwise unsuitable for the application.\nTherefore there is a need for new magnetic field configurations providing new magnetic field properties that can improve and extend the operation of existing magnetic field devices and potentially bring the benefits of magnetic field operation to new devices and applications heretofore served only by purely mechanical devices."}
{"text": "The present invention relates to dollies, or carts, adapted to transport large sheets of material across floors and the like.\nIt has long been a problem to move single large sheets of metal, glass, stone and ceramic about in manufacturing plants and in the field. As an example, a conference table may have a large stone or glass sheet top, for instance, a 5xc3x9710 foot rectanguloid or 6 foot diameter disk weighing 300 or more pounds. Typically, such a top rests on a detachable base, to facilitate transport. When such a large table top is be carried into an office building lobby, up an elevator, and through a finished office, care is needed to protect the top and things along the path. The same problems obtain when glaziers transport large sheets of glass any distance. As another example, in a fabrication shop, a sheet of steel or other material might need to be casually moved from one work station to another. Not only are such large sheet objects awkward to handle, but particularly in the field, it is important that the sheet be handled in a safe way which damages neither the sheet nor other things.\nA common way to handle a large sheet is to have at least two or more workers manually carry the sheet. Other times, a conventional flat dolly is employed, even though the sheet tends to be unstable and prone to come off such a type unit. Lashing the sheet to the dolly is awkward and inconvenient.\nPrior inventors have addressed the problem in various ways by designing special dollies. Many such dollies have xe2x80x9cAxe2x80x9d frames or other sloped structures mounted on their surface, to support the sheet. The angle at which the sheet is carried is fixed. The prior art dollies are often not suited for carrying round or oval sheets, as characterize many table tops. In general, the prior art A-frame type dollies are bulky and heavy; and, they are costly. They may be suited for certain factory settings. But, in the field it is inconvenient to transport such dollies due to their bulk and weight. The bottom line is that the dollies taught by the prior art tend to not be used because they are unsuitable for applications where a fragile environment, e.g., a finished office or living space, must be traversed.\nThus, there is a continuing need for a dolly which stably carries both rectangular and flat sheets, and which dolly is light, transportable, and economic to produce. And, experience shows that this may not be easy, because a large sheet mounted on a dolly can easily comprise an unstable combination.\nAn object of the invention is to provide means for transporting both large rectangular and large circular sheet objects. A further object of the invention is that such means be light in weight, convenient to use, easy to transport, and economic to fabricate.\nIn accordance with the invention, a sheet to be transported is placed edgewise on the base of a rectangular frame dolly base having four wheels, preferably swivel caster wheels. The frame of the dolly base is formed of two lengthwise running side beams and two opposing end members; and, it has a central rectangular opening. A sheet being transported leans against the vertical surfaces of a stanchion comprised of spindles and a horizontal top bar. The stanchion is mounted offset from the longitudinal centerplane of the base. As a rectangular sheet is being carried, it rests along a plane defined by the top surfaces of the front and rear members. When a circular sheet is transported it is cradled within the space between the front and rear members, and its lower edge is at an elevation less than said plane.\nFor stability during transport, the sheet center of gravity lies within the vertical projection of a no-contact zone, as preferably do the lower and upper parts of a sheet which contacts respectively the base and the stanchion. The no-contact zone is a region defined by the innermost portions of circles of rotation of swivel caster wheels where they touch the floor and swivel around their pivot axes. The inner surfaces of the stanchion which contacts the sheet lie along a vertical plane. That vertical plane is located within bounds of the no-contact zone, called Zone Z herein. Likewise, on the opposing side of the base, the inside of the frame hole of the base is also located within the no-contact zone vertical projection, to desirably position the bottoms of circular sheets placed on the dolly. More preferably, the stanchion and interior of the base frame are located within a sub-portion of the Zone Z, which is called Zone Y. Zone Y is defined by the innermost portions of the circles of rotation of the outer circumferences of the swivel wheels. In the preferred invention, the sheet lies against the stanchion at an angle with the horizontal of 75-90 degrees, preferably 75-85 degrees, most preferably 80-85 degrees.\nPreferably, there a retention bar sticking up from the base surface at said opposing side, within the no-contact zone bounds, to position the bottom edge of a rectangular sheet. The retention bar serves the same purpose as the interior of the base frame does for circular sheets; and, when the retention bar extends across the opening of the base frame, it is substitutional for the inner edge of the base frame in determining where a circular sheet bottom will rest. The top surfaces of the end members comprises a resilient frictional material, e.g., rubber underlain by wood. Since the sheet weight is concentrated at two spaced apart points, local deformation of the resilient surface helps stop sideways slipping of the sheet. More preferably, the stanchion, frame hole and retention bar components of the dolly are analogously positioned with respect to a smaller zone, called the no-wheel zone, which lies within the no-contact zone. The no-wheel zone is defined by the innermost circumferential points of the wheels themselves, as compared to where they contact the floor.\nIn further accord with the invention a second horizontal bar, which is vertically adjustable, runs between the spindles of the stanchion. This provides an adjustable support which is particularly useful for small circular sheets. Preferably, the stanchion is removable from the base, by having the spindles mounted in sockets, or clamped to the base. Alternatively, the stanchion is hinged to the base, so it can be folded flat along the top surface. These features facilitate transport of the dolly to the point of use.\nThe dolly of the invention is useful for stably transporting a variety of shapes of sheets and economic to construct. The dolly configuration minimizes chance for damage to the sheet edges or by bending fracture during use. And even while providing these advantages, the dolly is light weight, adapted to easy transport, and economically made.\nThe foregoing and other objects, features and advantages of the invention will become more apparent from the following description of the best mode of the invention and accompanying drawings."}
{"text": "This invention relates to the steering of beams of electromagnetic radiation, such as light beams, by relative translation of lens arrays in combination with phase shifters.\nCoherent beams of electromagnetic radiation are scanned for use in communication systems, radar, weapons, welding, supermarket label checking, and optical disc reading and writing. Very often, the transmitted beams are made up from a combination of plural individual beams.\nThe scanning function may be provided by gimballed, mechanically moveable mirrors, lenses or reflectors. However, the mass of such structures may impede the ability to scan in a random fashion, although repetitive scanning at high speeds may be possible. An article entitled \"Binary micro optics: an application to beam steering\", by Goltsos et al., published by Lincoln Laboratory in connection with the SPIE: OE LASE 89, 1052 (January 89) describes the relative translation of a pair of microlens arrays for beam steering. As described in the article, beam steering is accomplished by relative translation of a pair of microlens arrays cascaded in the path of an array of light beams. The translation of the microlens arrays is in a direction lateral to the beam direction, and the magnitude of the motion which is required for scanning is less than the diameter of the individual lens of the array.\nFIG. 1a illustrates a portion of a cascade of two microlens arrays. In FIG. 1a, a scanner designated generally as 10 includes a first microlens array 12 which includes individual lenses 14, 16, 18 and 20. Adjacent light beams illustrated as 22, 24, 26 and 28 fill the apertures of lenses 14, 16, 18 and 20, respectively. Lenses 14-20 cause the light beams to converge toward focal points (not illustrated). A second microlens array 32 includes diverging or defocussing lenses 34, 36, 38 and 40. Microlens array 32 is capable of translation relative to microlens array 12 in a direction of arrows 41. When the lenses of the microlens arrays 12 and 32 are registered, i.e., when the corresponding lenses are coaxial as illustrated in FIG. 1a, the output light beams, illustrated as 42, 44, 46 and 48, propagate parallel to the direction of propagation of incoming light beams 22, 24, 26 and 28, respectively.\nFIG. 1b illustrates as plots 52, 54, 56 and 58 the phase of the wave fronts associated with light beams 42, 44, 46 and 48, respectively, as a function of distance from an arbitrary reference point relative to the lens arrays. The spaces between plots 52, 54, 56 and 58 represent regions in which the light beams have a small amplitude. In FIG. 1b, plots 52, 54, 56 and 58 are, in effect, portions or continuations of the same straight dash-line 51 having the same phase. Other plots could be made at other distances from the lens arrays, with the phases increasing gradually with increasing distance from the lens arrays, and with the phases recurring if reduced by subtraction of multiples of 2.pi..\nAs illustrated in FIG. 1a, the output apertures of the lenses of array 32 are not filled. If the output apertures were filled, plots 52, 54, 56 and 58 of FIG. 1b would run together to create a continuous phase front representing a coherent beam of light, the direction of propagation of which is normal to the phase front.\nFIG. 1c illustrates a portion of scanner 10 of FIG. 1a, with lenses 36 and 38 of lens array 32 translated vertically upward (in the direction of arrow I) relative to corresponding lenses 16 and 18 of lens array 12, and with the input light beams 24 and 26 illustrated as not completely filling the input aperture to enable the beam paths to be clearly depicted. As illustrated, output beams 44 and 46 propagate in a direction different from that of the incoming beams, i.e. the beams have been scanned. FIG. 1d illustrates the phase of the wave fronts of beams 44 and 46. As illustrated in FIG. 1b, phase fronts 64 and 66 exhibit a slope, the normal to which defines the direction of propagation of the beam. As also illustrated in FIG. 1d, there is an offset, which is illustrated between arrows 50, which represents the offset between the phases of adjacent continuations of beams 44 and 46 of FIG. 1c. If this phase offset is zero or zero plus a multiple of 2.pi., the beams are in-phase for the illustrated direction of propagation, and a beam maximum occurs. In general, however, the phase offset will vary with the scanning direction, with the result that for some scanning directions the individual beams will be mutually out-of-phase with another beam, resulting in destructive interference. This in turn results in a far-field scanned radiation pattern which contains grating lobes or angles at which the radiated energy is high, and other angles at which the radiated energy is low. The result of translating a lens array in one direction is to gradually reduce the amplitude of one grating lobe, while the adjacent grating lobe becomes larger. The Goltsos et al. article suggests the use of a scanning mirror at the system input for fine or vernier beam steering. Such a scanning mirror has the disadvantages of a mechanical system referred to above, and in addition, causes the beams to enter the lenses of the lens array at an angle, which reduces the efficiency of the lens. This may be particularly important when two lens arrays are involved, because the entry at an angle occurs in both lens arrays, so the losses are cascaded. It is desirable to scan in a manner which allows the beam(s) to be directed at any angle, and not just at angles at which grating lobes occur.\nFIG. 2a illustrates a lens array similar to that of FIG. 1, with the lenses of the two arrays registered, and FIG. 2b illustrates the same arrangement with one of the arrays laterally offset by translation in the direction of arrow I. In FIG. 2, elements corresponding to those of FIG. 1 are designated by the same reference numerals. In FIG. 2a, circular input light beams 24 and 26 are centered on axes 6 and 8, respectively, and fill the apertures of converging lenses 16 and 18, pursuant to the Goltsos et al. suggestion. Lenses 16 and 18 focus the light to form converging beam portions 74 and 76, respectively, which come to a focus at a focus plane 99. From focus plane 99, diverging beam portions 84 and 86 propagate toward the input apertures of converging lenses 234 and 236, respectively. As illustrated, the spacings are such that beam portions 84 and 86 do not fill the apertures of lenses 234 and 236. Lenses 234 and 236 collimate the beams to produce parallel output beams 44 and 46, respectively, which are centered on axes 4 and 6, respectively.\nFIG. 2b illustrates the result of moving lens array 32 of FIG. 1a downward, in the direction of arrow I. As illustrated, light beams 84 and 86 intercept lenses 234 and 236 in a region in which the lens curvature causes output beams 44 and 46 to be deflected or scanned downward.\nFIG. 3a is identical in subject matter to FIG. 1c, and is included as a reference for comparison with FIG. 3b. In FIG. 3b, array 32, which includes diverging lenses 36 and 38, has been moved or translated upward in the direction of arrow I, thereby causing exit beams 44 and 46 to be deflected downward.\nBy comparison of FIGS. 2b and 3b, it is apparent that deflection of output beams in a given direction in accordance with the Goltsos et al. arrangement requires that the output lens array be moved in the direction of the desired deflection in the case of converging lens array, and in a direction opposite to the desired scanning direction for a diverging lens array.\nIn FIG. 4, an arrangement similar to that of FIG. 2 has had its output lens array 32 translated upward by an amount .DELTA. in an attempt to increase the scan angle. As illustrated, translation .DELTA. is sufficient to cause diverging light beam portion 84 to illuminate portions of both lenses 234 and 236. This may be viewed as a form of overfilling of the aperture of lens 234. As illustrated, output beam 244 is deflected or scanned upward by lens 234. That portion of beam 84 falling onto lens 236, however, is deflected downward. When overfilling of the aperture occurs in this manner, the far-field peak beam amplitude decreases, because the effective aperture decreases. Put another way, translation of the moving lens array in the direction of the arrow in any of FIGS. 2, 3 or 4 may result in a secondary portion of each beam (246) being directed away from the main scanned beam (244). The energy which goes into the secondary beam is not available for the main beam, and the secondary beam amounts to a scanning sidelobe which may not be desired. The undesirable effect of overfilling also occurs with the arrangement of FIG. 3. It would be advantageous to be able to translate the lenses to achieve additional scanning, with less loss of peak amplitude."}
{"text": "The present invention relates to an apparatus for imprinting checks, money orders, and other negotiable instruments.\nExisting apparatus for imprinting negotiable instruments, employ moveable independent type segments, which when punctuation indicia is desirable, that is, periods, commas, or foreign currency symbols, such apparatus requires several internal and external components to be modified, redesigned, and retooled to accommodate such punctuation indicia. Additionally, existing check writers or apparatus for imprinting negotiable instruments employ an independently sliding prefix plate which allows the elimination of a space or series of easily altered symbols to be imprinted in front of the dollar amount. This prefix plate commonly imprints THE SUM tightly in front of the dollar amount, thus minimizing the opportunity for criminal alteration of the dollar amount. However, due to the various dollar amount capacity check writers manufactured, which may include anywhere from between five (5) to thirteen (13) independent type segments, several different sizes and configurations of these prefix plates must be produced to fit the various capacity check writers produced. Accordingly, this variation in size and configuration of prefix plates causes confusion in the ordering of replacement prefix plates as well as operational problems due to the incorrectly sized and dimensioned prefix plate malfunctioning in the check writer. Also, the present cross-sectional configuration of the prefix plate has proven to create oscillation inconsistencies and poor imprint quality.\nAt present, available manual check writers or apparatus for imprinting negotiable instruments use one of two common inking methods. The first method includes use of an ink roller or ink pad device. This method requires the operator to periodically reink the device which produces inconsistent imprints and presents a messy inconvenience to the operator. The second method includes the use of an inked fabric ribbon. This inking method also requires the operator to periodically change the internal ribbon, an operation which has proven to be a difficult and messy procedure. For this reason, maximum security dye based inks are not commonly used in the manufacture of the ribbon fabric of such ink ribbon assemblies for check writing machines. Another problem encountered with the use of an ink ribbon assembly in check writers involves the self-reversing mechanism used with such ink ribbon assemblies. Because the self-reversing mechanism is a tension driven application, successful activation of the mechanism requires that proper tension be maintained at all times on both the ribbon and take-up spools, a result which is difficult to control and to maintain with such replacement ink ribbon assemblies."}
{"text": "Butanol is an important industrial chemical, useful as a fuel additive, as a feedstock chemical in the plastics industry, and as a foodgrade extractant in the food and flavor industry. Each year 10 to 12 billion pounds of butanol are produced by petrochemical means and the need for this commodity chemical will likely increase. 2-Butanone, also referred to as methyl ethyl ketone (MEK), is a widely used solvent and is the most important commercially produced ketone, after acetone. It is used as a solvent for paints, resins, and adhesives, as well as a selective extractant and activator of oxidative reactions.\nMethods for the chemical synthesis of 2-butanone are known, such as by dehydrogenation of 2-butanol, or in a process where liquid butane is catalytically oxidized giving 2-butanone and acetic acid (Ullmann's Encyclopedia of Industrial Chemistry, 6th edition, 2003, Wiley-VCHVerlag GmbH and Co., Weinheim, Germany, Vol. 5, pp. 727-732). 2-Butanone may also be converted chemically to 2-butanol by hydrogenation (Breen et al., J. or Catalysis 236: 270-281 (2005)). Methods for the chemical synthesis of 2-butanol are known, such as n-butene hydration (Ullmann's Encyclopedia of Industrial Chemistry, 6th edition, 2003, Wiley-VCHVerlag GmbH and Co., Weinheim, Germany, Vol. 5, pp. 716-719). These processes use starting materials derived from petrochemicals and are generally expensive, and are not environmentally friendly. The production of 2-butanone and 2-butanol from plant-derived raw materials would minimize greenhouse gas emissions and would represent an advance in the art.\nMethods for producing 2-butanol by biotransformation of other organic chemicals are also known. For example, Stampfer et al. (WO 03/078615) describe the production of secondary alcohols, such as 2-butanol, by the reduction of ketones which is catalyzed by an alcohol dehydrogenase enzyme obtained from Rhodococcus ruber. Similarly, Kojima et al. (EP 0645453) describe a method for preparing secondary alcohols, such as 2-butanol, by reduction of ketones which is catalyzed by a secondary alcohol dehydrogenase enzyme obtained from Candida parapsilosis. Additionally, Kuehnle et al. (EP 1149918) describe a process that produces both 1-butanol and 2-butanol by the oxidation of hydrocarbons by various strains of Rhodococcus ruber. The process favored 1-butanol production with a selectivity of 93.8%.\nThe production of 2-butanol by certain strains of Lactobacilli is also known (Speranza et. al. J. Agric. Food Chem. (1997) 45:3476-3480). The 2-butanol is produced by the transformation of meso-2,3-butanediol. The production of 2-butanol from acetolactate and acetoin by these Lactobacilli strains was also demonstrated. However, there have been no reports of a recombinant microorganism designed to produce 2-butanol.\nThere is a need, therefore, for environmentally responsible, cost-effective processes for the production of 2-butanol and 2-butanone. The present invention addresses this need through the discovery of recombinant microbial production hosts expressing 2-butanol and 2-butanone biosynthetic pathways."}
{"text": "Various devices/systems may include an embedded platform comprising an operating system, applications and various other files. For example, the various devices/systems may include portable ultrasound machines, Global Positioning System (GPS) devices, Automated Teller Machines (ATMs), devices that power large construction machinery and/or the like. The operating system, the applications and important user data may be configured into a run-time image on which the devices/systems boot up and operate. For instance, Windows Embedded Technology, including the Embedded Windows operating system, integrates an existing Information Technology (IT) infrastructure and enables various services, such as management, security, data synchronization with a network or another device/system, usage profiling, location services, advertising services, business intelligence and line-of-business applications, access to data regarding device capabilities and services and/or the like.\nThe Embedded Windows operating system includes an optional feature known as filters (e.g., write filters), which redirect a user's change made on the operating system to an overlay, instead of making the change on the hard disk. As result, the operating system is protected from malicious/accidental modifies and remains in a same state through a device/system reboot. The current design of write filters may employ a RAM-based or a disk-based overlay (e.g., hard disk) as the overlay type. A size limit for the overlay media is predefined, such that filter may access part of a physical RAM or hard disk without impacting the user's normal activities. As more changes are made to the operating system, the overlay increases in size (e.g., a number of bytes). Eventually, the overlay media exceeds the size limit, which causes a crash, an unexpected reboot and other deleterious effects."}
{"text": "Systems that track objects in an event, such as participants in an athletic event, are known. For example, U.S. Patent Application Publication No. 2008/0129825 to DeAngelis et al., which is incorporated herein by reference, discloses systems and methods to facilitate autonomous image capture and picture production. A location unit is attached to each tracked object (e.g., participants in an athletic event). An object tracking device receives location information from each location unit. A camera control device controls, based upon the location information, at least one motorized camera to capture image data of at least one tracked object.\nIt is also known to manually create video and still images of an event. For example, a video feed of an event (e.g., an American football game) is typically generated by highly trained camera persons and highly trained production staff who select camera shots and combine graphics into the video feed. Video images and/or still picture production can be partially or fully automated using systems and methods disclosed in U.S. Patent Application Publication No. 2008/0129825.\nIn many American football games, two ‘standard’ views are manually filmed using two digital video cameras; one on the sideline, and one in an end zone. These views are then manually ‘broken down’ by humans watching the videos, clipping them into plays, and identifying interesting attributes of each play. One of the most obvious attributes is simply who was on the field for each team at a given time. This is also one of the most difficult things to determine from the video since the resolution is not sufficient to clearly determine each of the players' numbers, thus making it difficult to identify all of the players."}
{"text": "Generally speaking instantaneous feed error between the Label Stop Sensor (LSS) and the Thermal Print Head's (TPH) burn line always varies depending upon the type of label, forces acting upon label and ambient conditions. The LSS is a positional sensor, identifying the edge or gap or black mark of the label.\nWithout instantaneous feed error correction the quality of the print registration would be challenged. Print registration is the accuracy of the position of the printed image on the label and effects print quality.\nThere are systems known in the art for determining label positions. For example US Publication 20130244872A1 discloses a thermal printer with an optical registration system especially for use with labels having fluorescent stripe patterns. However, no line feed correction calculation is provided for. Likewise, U.S. Pat. No. 8,029,083 discloses a label printer to determine the position of a label. However, the '083 reference makes no provision for the determination and correction of line feed error for the label.\nTherefore, a need exists for a system and method of determining the position of labels on the line feed of a label printer, determining the line feed error of the label, and correcting the line feed error before the burn line on the label printer."}
{"text": "Computing systems generally utilize non-volatile storage media, such as hard disk drives, to store information, and to ensure the availability of stored information across power cycles. In some instances, computing systems or applications executed by computing systems require that data be committed to a disk drive in a specific order, such that subsequent read requests to the disk drive result in retrieval of the expected data.\nHowever, the systems and protocols used to interact with storage devices such as disk drives may be unable to ensure such a specific order of writes. For example, multithreaded systems may allow multiple applications or processes to access a disk drive simultaneously or near-simultaneously, and in doing so, may reorder various operations performed by the applications or processes. Further, current disk drives or other non-volatile storage media can utilize volatile cache memory in an attempt to increase speed and efficiency of the drive. Use of cache memory, and protocols associated therewith, may result in reordering of operations prior to commitment to a non-volatile memory. Reordering of operations to the non-volatile memory may result in unexpected values being written to or read from the non-volatile memory by the computing system, and therefore result in errors.\nTo compensate for the out-of-order environment in which writes can be committed to the non-volatile memory, an application has the option to ensure that all of its previous writes have completed before completion of the next write by performing an application file operation flush, which causes all data of all previously completed writes to be forced out of the system and committed to the storage device's permanent media. However, implementation of a flush command on the storage device is generally time consuming, and can degrade performance. Consequently, many storage subsystems (or software/firmware/hardware associated therewith) decline to perform flush commands. Ignoring a program's (e.g., application's or operating system's) flush commands may increase performance relative to performing the commands, but results in potential data inconsistency.\nAn alternative technique for ensuring data consistency is to serialize all ordered writes throughout the host computing device and mark each one as a Forced Unit Access (FUA). Use of FUA writes ensures that each write is committed to non-volatile media before the command is completed. However, due to the delay in completion, FUA commands tend to be significantly slower than a typical write command. As a result, many storage subsystems (or software/firmware/hardware thereof) fail to respect a program's request to perform FUA writes."}
{"text": "This invention relates to an alcohol reclaiming circuit for use in an absorption refrigeration system for maintaining the system at a constant high level of operation.\nIt has long been recognized in the absorption refrigeration art that the addition of alcohol to the working fluids contained within the system greatly enhances the overall performance of the system. Although the exact mechanism involved is not fully understood, the results obtainable can be dramatically demonstrated in practice. Often times, an increase in performance in excess of thirty percent (30%) can be immediately realized by the introduction of a small amount of alcohol into the system. However, it has heretofore been difficult, if not impossible, to sustain this high level of operation for any appreciable period of time. Experience has shown that the gain in performance that is initially achieved continually diminishes with time until the gain is ultimately reduced to zero. It has heretofore been the common practice to add more alcohol to the system When performance drops below a predetermined level. This procedure, however, is not entirely satisfactory in that it requires continual monitoring of the system's performance, produces an undesirable buildup of alcohol within the system, and results in unwanted variations in overall machine performance."}
{"text": "1. Field of the Invention\nThe present invention relates to methods of forming a chemical casing in a well bore penetrating a weak unconsolidated zone or formation.\n2. Description of the Prior Art\nRotary drilling methods are commonly utilized in the drilling of oil and gas wells. That is, the well bore which extends from the surface into one or more subterranean oil and/or gas producing formations is drilled by a rotary drilling rig on the surface which rotates a drill bit attached to a string of drill pipe. The drill bit includes rotatable cutting surfaces so that when the drill bit is rotated by the drill string against subterranean strata under pressure a bore hole is produced.\nA drilling fluid is circulated downwardly through the drill string, through the drill bit and upwardly in the annulus between the walls of the well bore and the drill string. The drilling fluid functions to maintain hydrostatic pressure on formations penetrated by the well bore and to remove cuttings from the well bore. As the drilling fluid is circulated, a filter cake of solids from the drilling fluid forms on the walls of the well bore. The filter cake build-up is a result of initial fluid loss into permeable formations and zones penetrated by the well bore. The presence of the filter cake reduces additional fluid loss as the well is drilled.\nIn addition to removing cuttings from the well bore and forming filter cake on the well bore, the drilling fluid cools and lubricates the drill bit and exerts a hydrostatic pressure against the well bore walls to prevent blow-outs, i.e., to prevent pressurized formation fluids from flowing into the well bore when formations containing the pressurized fluids are penetrated. The hydrostatic pressure created by the drilling fluid in the well bore may fracture low mechanical strength formations penetrated by the well bore which allows drilling fluid to be lost into the formations. When this occurs, the drilling of the well bore must be stopped and remedial steps taken to seal the fractures which are time consuming and expensive.\nIn order to insure that fracturing of low mechanical strength formations penetrated by the well bore and other similar problems do not occur, it has heretofore been the practice to intermittently seal the well bore by cementing pipe referred to in the art as casing or liners in the well bore. The points in the well bore during its drilling at which the drilling is stopped and casing or liners are installed in the well bore are commonly referred to as “casing points”. Casing or a liner is placed in the well bore above each casing point and a sealing composition such as a hydraulic cement composition is pumped into the annular space between the walls of the well bore and the exterior surface of the casing or liner disposed therein. The hydraulic cement composition is permitted to set in the annulus thereby forming an annular sheath of hardened substantially impermeable cement therein. The cement sheath physically supports and positions the pipe in the well bore and bonds the pipe to the walls of the well bore whereby the undesirable migration of fluids between zones or formations penetrated by the well bore is prevented. This technique of cementing pipe in the well bore as the drilling progresses has a number of disadvantages including the time and expense incurred in placing and sealing the pipe as well as the reduction in the well diameter after each casing point. That is, the well diameter must be reduced below each casing point so that a smaller casing can be lowered through the previously placed casing and sealed in the well bore.\nAnother problem that occurs in the drilling and completion of well bores is that when the well bore is drilled into and through unconsolidated weak zones or formations formed of clays, shales, sand stone and the like, unconsolidated clay, shale and sand slough off the sides of the well bore which enlarges the well bore and often causes the drill bit and drill pipe to become stuck whereby drilling must be stopped and remedial steps taken.\nThus, there are needs for improved methods of drilling well bores whereby unconsolidated weak zones or formations are consolidated and the mechanical strength of the well bore is increased during drilling without the need to stop drilling for prolonged periods of time."}
{"text": "Recent video coding schemes such as HEVC (High Efficiency Video Coding), MVC (Multiview Video Coding) or SVC (Scalable Video Coding) support inter-picture prediction using previously coded reference pictures. The classification of reference pictures may be different depending on, e.g., their Picture Order Count (POC) distance from a target picture, view layer for SVC and view id for 3D video and MVC. Nevertheless, conventional coding scheme applies the same motion vector prediction for all types of reference pictures regardless of their view id, view layers and whether they are short or long term reference pictures. For example, performing POC-based scaling on a motion vector (MV) which points to a long term reference picture (LTRP) may result in the scaled motion vector having an extremely large or small magnitude. In such a case, the accuracy and efficiency of the motion vector prediction process becomes suboptimal.\nHEVC supports spatial motion vector prediction and temporal motion vector prediction. In spatial motion vector prediction, the motion vector of a target prediction unit (PU) is predicted using a motion vector of a previously coded neighbouring PU. Both the target PU and the neighbouring PU are located within a current target picture. In temporal motion vector prediction, a motion vector of a target prediction unit (PU) is predicted using a motion vector of a collocated block. The collocated block is located within a previously coded collocated picture and the collocated block is coded using a motion vector pointing to a reference picture (which may be referred to as a collocated reference picture). The term collocated generally indicates that the coordinates of the collocated PU within the collocated picture are the same as the coordinates of the target PU within the current target picture. However, due to variable coding unit and prediction unit sizes in HEVC, the current PU and collocated PU may not be perfectly aligned (i.e. their coordinates may not be exactly the same), and a predetermined selection scheme is used for selecting the collocated PU.\nA motion vector predictor (MVP) may be obtained by scaling the motion vector of a neighbouring PU or a collocated PU based on certain characteristics of the motion vector such as its temporal distance (i.e., picture order count (POC) value difference) between the target picture and its corresponding reference picture. For example, the motion vector of a collocated PU may be scaled according to the POC distance to produce a temporal MVP for the current PU according to the following equation:MVP=(tb/td)*nmv \nwhere:                MVP=temporal motion vector predictor derived from the motion vector of collocated block/PU;        nmv=motion vector of the collocated block/PU;        tb=signed POC distance/difference from the current picture to the reference picture referred by the target block/PU;        td=signed POC distance/difference from the collocated picture to the reference picture referred by the collocated block/PU.        \nGenerally, in spatial prediction, for deriving “td”, the target picture is the current picture and its reference picture is the collocated reference picture. In temporal prediction, for deriving “tb”, the target picture is the collocated picture and its reference picture is the collocated reference picture. For deriving “tb”, target picture is the current picture and its reference picture is the current reference (i.e., referred by the target PU, either from RefList0 or RefList1) for both spatial and temporal prediction."}
{"text": "The coal to ethylene glycol technology mainly includes three steps: a first step, catalytically elimination small amount of hydrogen gas in CO separated by pressure swing adsorption from coal-derived syngas; a second step, vapour-phase CO oxidative coupling to oxalate; and a third step, hydrogenation of oxalate to ethylene glycol. Among them, vapour-phase CO oxidative coupling to oxalate is the key step to realize the conversion of inorganic CO to organic chemical oxalate in coal to ethylene glycol. Oxalates are important organic chemical materials useful for preparing ethylene glycol, oxalic acid, oxalyl chloride, oxalic amide, some medicaments, dyes, and intermediates of the solvents. Currently, oxalates are produced by the process of esterifying dehydration from oxalic acid and alcohol employing toluene or benzene as the dehydrator. This production process has a high process cost, a large energy consumption, serious pollution, and unreasonable material utilization. In 1960s, American, Fenton reported a method for producing oxalates from CO and alcohol by direct coupling, opening a new synthesis routine of oxalates by C1 chemistry. However, this reaction needs to be carried out under the condition of high pressure. Patent JP 8242.656 published patent disclosure reports a process for synthesizing dimethyl oxalate from CO and methyl nitrite under normal pressure employing supported Pd catalyst. The process has good economic efficiency, mild reaction conditions, low energy consumption, no discharge of waste gas, waste liquid and waste solid, and good product quality. The space-time yield of catalyst reported by this patent is 432 g·L−1·h−1, which does not decrease after 480 hours of continuous reaction. However, the amount of the noble metal used is large, resulting in a high cost of the catalyst, and the space-time yield is low. After that, lots of patents subsequently report adding promoters such as Zr (CN95116136.9), Ce (CN02111624.5), Ti (CN200710061392.2), La (CN200810114383.X), Ir (CN200810035248.6), Ni (CN200910307543.7), Cu (CN200910060087.0), MOx (CN200910061854.X), and the like in the catalyst to improve the space-time yield of oxalate. After adding promoter, the space-time yield of oxalate increased to a certain degree, but the amount of the active component Pd is still relatively high. Currently, the actual loading amount of Pd in the industrial plant is about 2%. In addition, all the catalysts reported in the patents and literatures are produced by traditional wet impregnation process, involving steps of immersing, drying, calcinating, reducing at high temperature, and the like. It takes long duration time and high energy consumption. Moreover, it is impossible to conduct the exact control to the size and exposed crystal facet of the Pd particles. However, the size and exposed crystal facet of the Pd particles are two very important parameters influencing the catalytic activity. Therefore, it is especially economically valuable to develop a simple and fast process for producing nanocatalysts with low noble metal loading amount, controllable size and exposed crystal facet as well as high activity for vapour-phase CO oxidative coupling to oxalate."}
{"text": "Without limiting the general scope of the invention, its background is described in connection with computer graphics, as one example only.\nIn computer systems, a host computer can be programmed to perform general purpose tasks including graphics routines. Greater speed and additional features are often desirable, and so a graphics processor is added to supplement the capabilities of the host computer.\nThe graphics processor is also called a graphics system processor (GSP), examples of which are the Texas Instruments TMS34010 and TMS34020 GSPs. The addition of a graphics processor makes the computer system a multiprocessor system which can benefit from advances in the art of multiprocessor technology. Furthermore, several different kinds of memory such as ROM, DRAM (dynamic random access memory) and VRAM (video RAM) are useful with computers that have graphics capability, and are desirably accommodated.\nIn computer graphics systems the low cost of dynamic random access memories (DRAM and VRAM) has made it economical to provide a bit map or pixel map memory for the system. In such a bit map or pixel map memory a color code is stored in a memory location corresponding to each pixel to be displayed. A video system is provided which recalls the color codes for each pixel and generates a raster scan video signal corresponding to the recalled color codes. Thus, the data stored in the memory determines the display by determining the color generated for each pixel (picture element) of the display.\nThe desirability of a natural looking display and the minimization of memory are conflicting. In order to have a natural looking display it is generally desirable to have a large number of available colors. This implies a large number of bits for each pixel in order to specify the particular color from among a large number of possibilities. However, the provision of a large number of bits per pixel calls for a large amount of memory for storage. Since a number of bits must be provided for each pixel in the display, even a modest sized display would require a large memory. Thus, it is advantageous to provide some method to reduce the amount of memory needed to store the display while retaining the capability of choosing among a large number of colors.\nThe provision of a circuit called a color palette enables a compromise between these conflicting requirements. The color palette stores color data words that are longer in bit length than color codes that are stored in the pixel map memory instead of the actual color data words themselves. The color data words can specify colors to be displayed in a form that is ready for digital-to-analog conversion directly from the palette. The color codes stored in the memory for each pixel have a limited number of bits, thereby reducing the memory requirements. The color codes are employed to select one of a number of color registers or palette locations. Thus, the color codes do not themselves define colors but instead identify a selected palette location. These color registers or palette locations each store color data words which are longer than the color codes in the pixel map memory. The number of such color registers or palette locations provided in the color palette is equal to the number of selections provided by the color codes. For example a four-bit color code can be used to select 2-to-the-n or sixteen palette locations. The color data words can be redefined in the palette from frame to frame to provide many more colors in an ongoing sequence of frames than are present in any one frame.\nDue to the advantages of the color palette devices, systems and methods, any improvements in their implementation are advantageous in computer color graphics technology. Indeed, any improvements in applicable circuits are desirable so that graphics and other computer and electronic systems can be made faster, more reliable and more convenient in commercial applications."}
{"text": "1. Field of the Invention\nThe present invention relates to a fixed wireless terminal (FWT) which is constructed by a mobile phone, and more particularly to an apparatus and a method for limiting a communication range of a fixed wireless terminal, which interrupt communication through a mobile phone when the mobile phone is not within a predetermined range from a fixed device.\n2. Description of the Prior Art\nIn the conventional art, in order to use a telephone in a mountainous region or an island, it is necessary to install a telephone line, which requires a significant expense.\nIn this consideration, recently, telecommunication through a fixed wireless terminal (FWT) has been disclosed in the art. The fixed wireless terminal is fixed indoors to serve only as a fixed device.\nThe FWT allows a mobile phone to be used in a room as if it were a wired telephone at a place where a telephone line is not installed or in a region where a call charge when using the FWT is cheaper than that when using a wired telephone. The FWT is one type of communication services, which is suggested by a communication service provider.\nIn the FWT, as shown in FIG. 1, a fixed device 100 can communicate with a base station 150 using an antenna 120 of a mobile phone. Power is supplied to the fixed device 100 using an adapter 110. The fixed device 100 and a handset 140 are connected to each other by a cord 130 to allow the FWT to perform the same function as a wired telephone.\nFIG. 2 illustrates another conventional FWT. In this type of FWT, a battery of a mobile phone 130 is removed, and the mobile phone 130 is connected with a fixed device 100 by a cord 120. The fixed device 100 is connected with a power supply device 110 so that power can be supplied to the mobile phone 130 through the fixed device 100 and the mobile phone 130 can directly communicate with a base station 140 to perform the same function as a wired telephone.\nHowever, in these types of FWTs, since communication is made possible only at a place where a fixed device is installed, mobility of the terminal in a room cannot but be limited to a certain extent. That is to say, it is possible to make and receive telephone calls only at a place where the fixed device is installed.\nTherefore, because the FWT cannot be used as a cordless telephone and only can be used as if it were a wired telephone, it is inconvenient to use the FWT."}
{"text": "The present invention relates to fishing, and more particularly, to devices used to hold fishing poles or rods that free up a fisherman\"\"s hands for accomplishing other tasks.\nWhile fishing with a fishing pole, it is frequently necessary for a fisherman to use both hands for threading a line and attaching lures, hooks, leaders and other fishing tackle to the line. Other tasks requiring the use of both hands include baiting, cleaning fouled hooks and removing fish that have been caught. Freeing both hands from the fishing pole normally means that the fisherman must let go of the fishing pole, running the risk of losing it overboard or fouling the reel or line. When fishing from the shore, fishermen often lean the pole against a chair or pier railing, and the pole often falls down. On a vessel, fishermen often lean the pole over the gunwales, only to have it fall down as the vessel pitches from side to side. Besides the inconvenience of picking up a fishing pole, the fishing pole can be damaged when it falls. Furthermore, the fishing pole\"\"s line can become tangled when the fishing pole falls over, especially when several adjacent fishing poles fall together into a pile.\nOn frequent occasions fishermen prefer to relax and fish xe2x80x9chands free.xe2x80x9d Some lakes allow fishing with a two pole license, Fishermen would therefore like to be able to take their hands off of the handle of the fishing pole and let the pole rest in an upwardly inclined orientation so they can attend to other tasks, like eating a snack or opening a canned beverage. In the absence of auxiliary equipment for holding the fishing pole, this requires that the pole be leaned against a railing or an object, or laid on the ground or vessel floor. Sometimes a fisherman will awkwardly try to hold the pole between his or her legs.\nA number of fishing pole holders have been designed to alleviate the foregoing problems. Many have been patented in the United States beginning in the nineteenth century. One type of fishing pole holder that is used near the shoreline consists of a cup or other fishing pole handle receptacle attached to a spike which is driven into the dirt, mud or sand at an inclined angle. See, for example, U.S. Pat. No. 4,455,779 of Cosic granted Jun. 26, 1984. These spike-mounted fishing pole holders are often hard to drive into rocky soil, are not easily moved and are unsuitable for use on a pier or on a vessel. Many fishing pole holders in the form of free-standing frames have been developed for use on piers and shorelines, but they are bulky and immobile. See, for example, U.S. Pat. No. 5,533,295 of Hochberger granted Jul. 9, 1996. Numerous fishing pole holders have been developed in the form of brackets which attach to chairs, vessel gunwales and railings. See, for example, U.S. Pat. No. 5,325,620 of Reed et al. granted Jul. 5, 1994, U.S. Pat. No. 4,017,050 of Rosenau granted Apr. 12, 1977, and U.S. Pat. No. 4,682,438 of Arrow granted Jul. 28, 1987. They are often unduly complex and immobile. They also rust and are difficult to adjust. In many cases, they require permanent attachment with screws and the like.\nBody-mounted fishing pole holders have been developed in an attempt to overcome the foregoing drawbacks. A body-mounted fishing pole holder is particularly attractive to fishermen who like to walk along the bank of a lake or stream so they can cast into a deep spot. Such fishermen often end up casting from rocks or steeply inclined bank areas that make it very difficult to bait their fishing pole and remove fish once caught because the pole cannot be easily stood on its end or laid down. Body-mounted fishing pole holders that have heretofore been developed have not experienced widespread adoption and use because they have not been designed to afford maximum comfort and convenience, they have been too expensive to manufacture and sell at low cost and/or they have lacked the required durability for the harsh conditions encountered by fishermen. Many of these devices have consisted of complicated harnesses. See, for example, U.S. Pat. No. 4,858,354 of Butts granted Aug. 22, 1989, U.S. Pat. No. 3,282,482 of Scharsu granted Nov. 1, 1966, and U.S. Pat. No. 5,520,312 of Maddox granted May 28, 1996. Others have consisted of handle receptacles mounted to a waist belt, which devices unduly restrict motion, are ill suited for sitting, do not adequately stabilize the fishing pole and place the tip of the pole too far away from the fisherman\"\"s hands. See, for example, U.S. Pat. No. 1,174,319 of Hipwood granted Mar. 7, 1916, U.S. Pat. No. 3,874,573 of Fruscella et al. granted Apr. 1, 1975, U.S. Pat. No. 4,569,466 of Webber granted Feb. 11, 1986 and U.S. Pat. No. 5,386,932 of Gross granted Feb. 7, 1995.\nU.S. Pat. No. 5,511,336 of Bishop granted Apr. 30, 1996 discloses a fishing pole holder having a hip plate made of thin, flexible plastic and a cylindrical tube attached to the hip plate. The hip plate is worn on a fisherman\"\"s belt and the handle of the fishing pole is inserted into the tube. This holder places the tip of the fishing pole too far away from the fisherman\"\"s hands. It is not comfortable to wear, and it does not allow the fisherman to easily sit with the pole mounted therein. Due to its molded plastic construction, a significant investment in tooling would be required to manufacture the patented fishing pole holder of Bishop.\nU.S. Pat. No. 5,956,883 of Krouth et al. granted Sep. 28, 1999 discloses a fishing pole holder comprising a concave, convex frame attached to the fisherman\"\"s lower leg. The frame has a boss with a polygonal pilot hole. A pole holding assembly including a swivel mechanism has a polygonal pin which is inserted into the pilot hole to allow for angular adjustment. The patented fishing pole holder of Krouth et al. is too complex and expensive to manufacture, and not readily adapted to support various pole sizes, particularly larger poles. In addition, all of the load of the fishing pole is transferred to a small region adjacent the boss which can lead to uncomfortable chafing against the fisherman\"\"s lower leg. The pole holding assembly of Krouth et al. projects too far outward from the fisherman\"\"s leg and can snag on branches and boat hardware.\nU.S. Pat. No. 6,003,746 of Richardson granted Dec. 21, 1999 discloses a body-mounted, strap-on fishing pole holder including a stabilized vertical frame upon which is secured a rotatable rod receiving element which is adjustable to any desired vertical angle. Stabilization is achieved by the use of rigid U-shaped members that are attached at vertically spaced locations on the vertical frame and surround the waist and thigh of the fisherman. Belts are connected to the ends of the U-shaped members for encircling the corresponding body portion. The patented fishing pole holder of Richardson is overly complex and subject to mechanical failures. It also places the tip of the fishing pole too far away from the fisherman\"\"s hands.\nU.S. Pat. No. 6,138,976 of Fahringer, Sr. granted Oct. 21, 2000 discloses a fishing pole holder consisting of straps and snaps that connect the upper portion of the pole to the fisherman\"\"s wrist and the lower handle portion of the fishing pole to the fisherman\"\"s ankle. This design is not suited for sitting, and it is tedious to connect and disconnect the pole from the wrist and ankle each time the fisherman wants to bait, cast and unhook fish. The fishing pole holder of Fahringer, Sr. is not adapted to allow hands free fishing once the line has been cast.\nU.S. Pat. No. 1,761,497 of Smith granted Jun. 3, 1930 discloses a fishing pole holder held to the fisherman\"\"s thigh with a pair of adjustable straps. The holder comprises a flat metal plate with a lower socket device in which the butt-end of the pole is seated and an upper spring clasp device including a pair of resilient branched arms. A pair of semi-circular jaws are pivotally connected to the arms. The jaws may be closed about the fishing pole by pressing the rod against the pivot joint between the jaws. The patented fishing pole holder of Smith is overly complex from a mechanical standpoint, subject to rusting and breakage, and ill suited for holding a variety of fishing poles, particularly larger ones. The Smith holder does not conveniently place the tip of the pole in the best location and its metal plate would be very uncomfortable pressing against the fisherman\"\"s thigh.\nU.S. Pat. No. 6,029,872 of Ellington granted Feb. 29, 2000 discloses a plastic fishing pole holder with an upper end that clips over the waistband of a swimsuit, a lower end held by a strap to the fisherman\"\"s thigh, and an intermediate flared and slotted tubular rod holder. This device would be expensive to mold, uncomfortable to wear, and would not support the weight of larger fishing poles. In addition, the patented fishing pole holder of Ellington does not conveniently position the upper free portion of the fishing pole.\nIt can be safely said that there has been a long felt need for an improved, versatile body-mounted fishing pole holder that is convenient, comfortable, inexpensive, lightweight and durable.\nIt is therefore object of the present invention to provide an improved body-mounted fishing pole holder that allows for more convenient stowage and retrieval of a fishing pole.\nIt is another object of the present invention to provide an improved body-mounted fishing pole holder that will place the tip of the pole in a more convenient location for baiting and fish removal.\nIt is still another object of the present invention to provide an improved body-mounted fishing pole holder that is adapted for either standing or sitting.\nIt is yet another object of the present invention to provide an improved body-mounted fishing pole holder that is inexpensive to manufacture yet very durable.\nIt is a further object of the present invention to provide an improved body-mounted fishing pole holder that is easy to don and remove and comfortable to wear.\nIt is yet another object of the present invention to provide an improved body-mounted fishing pole holder that is lightweight.\nIt is still a further object of the present invention to provide an improved body-mounted fishing pole holder that can also provide convenient storage for accessories such as a knife and fisherman\"\"s pliers.\nIn accordance with the present invention a fishing pole holder comprises a holster sized to overlie an outer lower leg portion of a fisherman. The holster has a receptacle sized and configured for removably receiving a lower handle portion of a fishing pole. The receptacle provides the sole support for the fishing pole without requiring any auxiliary holder connected to any other part of the fishing pole or the use of a hand or an arm of the fisherman. Straps removably secure the holster to the lower leg portion of the fisherman. The holster is configured for inclining the lower handle portion of the fishing pole in a forward direction away from the outer lower leg portion of the fisherman so that an upper free portion of the fishing pole extends in a generally upright orientation. The upper free portion of the pole is therefore inclined in the forward direction away from the fisherman\"\"s torso to position the tip of the pole for convenient two-handed baiting of a hook attached to a fishing line of the fishing pole and two-handed removal of a fish from the hook.\nAccording to the preferred embodiment of the present invention, a fishing pole holder includes a holster sewn from durable waterproof synthetic woven fabric to provide a primary pocket. The holster is held to the outer side of a fisherman\"\"s lower leg by upper and lower fabric straps stitched to the holster and wrapped around the lower leg. The straps are releasably secured with adjustable mating VELCRO(copyright) barb and loop fabric sections. A plastic tube is supported in the primary pocket of the holster. The tube is dimensioned to slidingly receive and removably the handle of a fishing pole and provide its sole support. The tube is forwardly inclined relative to the fisherman\"\"s lower leg so that the upper free portion of the fishing pole extends in a forward direction away from the torso of the fisherman. A neoprene spacer block is secured to an upper end of the holster on its rear side for also inclining the tube laterally so that the upper free portion of the fishing pole also extends in a sideways direction away from the torso of the fisherman. The tip of the fishing pole is therefore positioned in front, and to the right, of the fisherman for convenient baiting of the hook and removal of fish that have been caught. The fishing pole holder is versatile in that it can be comfortably used in either a standing or sitting position. The holster is fabricated with a secondary pocket that forms a fishing knife sheath and a strap loop that holds a pair of fisherman\"\"s pliers.\nThe present invention also provides a method of fishing which involves the steps of removably securing a holster around the fisherman\"\"s lower leg, the holster having a forwardly inclined receptacle for removably receiving and supporting the lower handle portion of a fishing pole. The upper free portion of the fishing pole is inclined forwardly away from the fisherman\"\"s torso. By inserting the pole into the receptacle, the fisherman can easily accomplish two-handed baiting of the hook and removal of fish therefrom between casting and reeling in catches."}
{"text": "1. Field\nThe present disclosure relates to an FET (field effect transistor) characteristic measurement method in which a pulse voltage output from a pulse generator is applied to the gate of an FET in order to measure drain current flowing through the FET. More specifically, the present disclosure relates to a method for accurately measuring a voltage applied to the drain of the FET and the drain current.\n2. Description of the Related Art\nConventionally, the IV (current-voltage) characteristics of an FET are determined by applying a predetermined DC voltage to the gate of the FET while a predetermined bias voltage is applied to the drain of the FET.\nHowever, the known approach in which the DC voltage is applied to the gate has the following problems.\nWhen the known measurement approach is used for measuring the IV characteristics of SOI (silicon on insulator) MOSFETs, strained-silicon MOSFETs, or the like, reliable measurement results of the IV characteristics may not be obtained due to a self-heating phenomenon of the FETs.\nWhen the known measurement approach is used for measuring the IV characteristics of MOSFETs using high-k (high dielectric constant) gate insulators, reliable measurement results of the IV characteristics may not be obtained since electrons are trapped in defects in the insulator films, thus reducing the drain current driving force.\nAccordingly, it has been proposed to apply a measuring method in which a short-duration pulse is applied to the gate of the FET (e.g., refer to K. A. Jenkins and J. Y-C. Sun, IEEE Electron Device Letters, Vol. 16, No. 4, April 1995, pp. 145 to 147). With such a measuring method, measurement results that are not affected by the self-heating can be obtained for MOSFETs employing SOI and strained silicon. Furthermore, with respect to MOSFETs using high-k gate insulators, IV characteristics that are close to these under actual operational conditions can be obtained (since microprocessors using the MOSFETs have internal circuits driven by pulses, not DC), without being affected by, for example, a reduction in a drain-current driving force caused by electron trapping in insulator defects.\nWhen the measurement method using short-duration pulses is employed, a so-called bias tee is used. An DC (direct current) input of the bias tee is connected to a DC voltage source, a bias output of the bias tee is connected to the drain of an FET, and an AC output of the bias tee is connected to an input of a measuring apparatus, such as an oscilloscope. When the short-duration pulse is input to the gate of the FET, drain current generated in a pulsed manner is input to the measuring apparatus via the AC output of the bias tee and the current is then converted by the input impedance of the measuring apparatus into a voltage.\nFrom a high-frequency point of view, the input impedance of the measuring apparatus acts as a shunt resistor interposed between the drain of the FET and the DC voltage source. Thus, when the drain current is generated, the drain bias voltage of the FET decreases because of a voltage drop caused by the input impedance. In a measuring system including the bias tee, the FET, and the measuring apparatus, a drop in the drain bias voltage can also occur because of impedance due to other elements, such as cables.\nHowever, in the known measuring method, the value of the drain bias voltage corresponding to a measured drain current is regarded as the value of a voltage output from the DC voltage source during the measurement of the drain current (i.e., is regarded as the value of the bias output voltage of the bias tee), that is, a voltage drop caused by the input impedance of the measuring apparatus and so on is not taken into account. Therefore, the measured drain current is not based on an actual bias voltage of the drain terminal, thus causing a problem in that accurate measurement results of IV characteristics cannot be obtained.\nIn addition, in the known measuring method, the drain current may not be accurately detected because of current leaking from the bias tee and so on, thus making it difficult to obtain accurate measurements results of IV characteristics."}
{"text": "1. Field of the Invention\nThe present invention relates to a device for adjusting an ink supply gap for an ink fountain apparatus of a printing press, and more particularly to a device for adjusting an ink supply gap adapted to move blade segments toward or away from an ink fountain roller so as to adjust the width of the ink supply gap.\n2. Description of the Related Art\nIn an ink fountain apparatus for an inking arrangement of a printing press, ink is stored in an elongated space of a substantially V-shaped cross section, which is formed by the peripheral surface of an ink fountain roller, the bottom portion of the ink fountain apparatus having its forward edge located in proximity to the peripheral surface of the ink fountain roller with an appropriate gap therebetween, and a pair of side plates, each disposed at one lateral end of the bottom portion. When the ink fountain roller is rotated, the ink stored in the elongated space is withdrawn through an ink supply port, i.e., a gap between the peripheral surface of the ink fountain roller and the forward edge of the bottom portion of the ink fountain apparatus.\nTo control the amount of ink withdrawn through the ink supply port, i.e., the supply of ink, the ink fountain apparatus has a blade for adjusting ink supply located at its bottom portion and a device for adjusting an ink supply gap. The device for adjusting the ink supply gap causes the blade to advance or retract so that the forward edge of the blade moves toward or away from the peripheral surface of the ink fountain roller, thereby adjusting the opening of the ink supply port.\nSuch conventional devices for adjusting an ink supply gap are disclosed in, for example, Japanese Patent Application Laid-Open (kokai) Nos. 7-246699 and 8-230161.\nAccording to \"Blade Adjusting Device for Ink Fountain of Printing Press\" disclosed in Japanese Patent Application Laid-Open (kokai) No. 7-246699, in order to adjust the width of a gap between an ink fountain roller provided in an ink fountain of a printing press and each of a plurality of blade segments, a gap adjust device is provided for each of the blade segments. The gap adjustment device includes pushing means for continuously urging a blade segment toward its base end by means of a compression coil spring; a push rod disposed in contact with the base end of the blade segment and adapted to move the blade segment; a front cam whose support shaft is disposed in parallel with the moving direction of the blade segment and whose face moves in parallel with the moving direction of the blade segment and is in contact with the base end of the push rod; position adjustment means for moving the support shaft of the front cam axially so as to adjust the axial position of the front cam; and drive means for rotating the front cam.\nIn order to adjust the initial position of each blade segment, the corresponding drive means is operated to rotate the corresponding front cam until the blade segment advances to a foremost position closest to the ink fountain roller. Then, the position adjustment means is operated to axially move the support shaft of the front cam, thereby adjusting the axial position of the front cam so as to bring the forward edge of the blade segment into contact with the ink fountain roller. In this state, the position adjustment means is locked to thereby set the blade segment to its initial position.\nDuring printing, according to the amount of ink required at a widthwise position of a printing surface, the gap between a corresponding blade segment and the ink fountain roller is adjusted in the following manner. The drive means is operated to rotate the front cam in such a direction that the blade segment moves away from the ink fountain roller. Since the pushing means pushes the blade segment toward the front cam, the forward edge of the blade segment moves from the position of contact with the ink fountain roller to an appropriate position located away from the ink fountain roller.\nAccording to \"Regulating Device for Opening of Ink Key\" disclosed in Japanese Patent Application Laid-Open (kokai) No. 8-230161, in order to adjust a gap between an ink fountain roller provided in an ink fountain of a printing press and each of a plurality of ink keys (corresponding to blade segments in the present invention), a gap adjustment device is provided for each ink key. The gap adjustment device includes a first hydraulic cylinder equipped with a piston connected to an ink key; a second hydraulic cylinder connected to the first hydraulic cylinder via an oil line and serving as drive means; and a third hydraulic cylinder provided in an oil line branching from the oil line which connects the first hydraulic cylinder and the second hydraulic cylinder and serves as adjustment means.\nIn order to move the ink key to its initial position, where the upper forward edge of the ink key is in contact with the peripheral surface of the ink fountain roller, the second hydraulic cylinder serving as drive means is locked, and the piston of the third hydraulic cylinder serving as adjustment means is moved to thereby move the piston of the first hydraulic cylinder. As a result, the ink key connected to the piston of the first hydraulic cylinder is moved accordingly.\nThe third hydraulic cylinder serving as adjustment means is locked, and the piston of the second hydraulic cylinder serving as drive means is moved by a drive device to thereby move the piston of the first hydraulic cylinder by means of oil contained in the oil line. As a result, the ink key connected to the piston of the first hydraulic cylinder is moved accordingly.\nNotably, the cross sectional area of the pistons increases in sequence of the piston of the third hydraulic cylinder, that of the second hydraulic cylinder, and that of the first hydraulic cylinder. According to Pascal's principle, the stroke of the piston of the first hydraulic cylinder connected to the ink key becomes very small as compared to that of the second hydraulic cylinder serving as drive means and that of the third hydraulic cylinder serving as adjustment means. Accordingly, the movement of the ink key can be finely adjusted.\nThe above-described conventional gap adjustment devices involve drawbacks.\nSpecifically, in \"Blade Adjusting Device for Ink Fountain of Printing Press\" disclosed in Japanese Patent Application Laid-Open No. 7-246699, a plurality of blade segments are each provided with a motor and a plurality of gears which constitute the drive means; shafts, a cam, and a push rod which constitute the position adjustment means; and a compression coil spring serving as the pushing means. Thus, a large number of parts are used, and the mechanism of the device is complex. As a result, the frequency of failure is high, and maintenance and inspection are time consuming.\nAccuracy in adjusting the gap between the ink fountain roller and a blade segment cannot be improved due to the backlash between gears and the influence of machining accuracy of the cam face.\nIn \"Regulating Device for Opening of Ink Key\" disclosed in Japanese Patent Application Laid-Open No. 8-230161, each of a plurality of ink keys must be provided with hydraulic piping; thus, complex piping is involved. Further, each ink key must be provided with a plurality of hydraulic cylinders, including a cylinder for moving the ink key, a cylinder for driving the ink key, and a cylinder for adjusting the ink key. Thus, the device assumes a large-scaled configuration, resulting in an increased manufacturing cost. Also, use of hydraulic oil involves a potential for oil leakage."}
{"text": "The present invention is a new low-cost method of boosting the pressure of a gas from a lower level to a higher level, for any purpose, one of which may be to inject the gas into a higher-pressure receptacle for purposes of sale, or to carry out some mechanical function. More particularly, it relates to such method and system especially adapted to the economics of the equipment used to increase the pressure (to pressurize or compress) the gas from a marginally-profitable gas well to a level acceptable for insertion into a gas sales line.\nThe present invention relates to a new low-cost method of compressing natural gas from a gas well when the reservoir pressure of the producing zone has declined below the pressure of the gas sales line. The present invention takes low-cost standard oil-field equipment and combines them into a unit that will compress gas more economically than a conventional compressor. The present invention relates, specifically, to the use of a flexible bladder inside a steel vessel, to receive and temporarily store natural gas from a gas well. In order to increase the pressure of the temporarily-stored gas inside the bladder, a hydraulic fluid is pumped into the annulus between the outer walls of the bladder and the inner walls of the steel vessel. With continued pumping, the pressure of the hydraulic fluid will exceed the gas pressure inside the bladder and the bladder collapses in size which results in the gas inside the reduced-size container (the bladder) being elevated to a higher pressure. The higher-pressured gas is then transferred to a natural gas pipeline for sale, or to some other mechanical process served by the higher-pressured gas. In the absence of an indirect method of pressurizing the gas, a more-expensive mechanical compressor would have to be used which would increase the cost of recovering the natural gas and result in natural gas reserves being abandoned which could be recovered by using the less-costly indirect compressor.\nThe present invention is specifically-related to the recovery of economic gas reserves when well-head pressures decline below the minimum suction pressure for conventional mechanical compressors, however, other uses of the equipment may be to collect small amounts of gas produced with oil for which there is no gas pipeline connection, to pressurize it into storage and further pressurize it for disposition by truck transport to a pipeline. Other uses may be to pressurize low-pressure gas from a local utility pipeline (approx. 1 psig) to 1000-1500 psi for compressed natural gas (CNG) fuel in isolated areas.\nThe flexible bladder is a one-piece cylinder-liner, when filled with gas from a gas well, or other source of low-pressure gas, will inflate substantially to the interior walls of the steel cylinder. The steel cylinder is sized to accommodate a number of fill-empty cycles during a 24 hour period such that the total volume of gas processed in a 24 hour period is substantial and the revenues generated from the sale of such gas will return the investment therein in a short time. The bladder is made of rubberized nylon, or if by choice, some other member of the elastomer family of synthetic rubbers, compatible with natural gas, and certain other gases, and a hydraulic fluid composed of fresh water -antifreeze mix, or mineral hydraulic oil, sealed at one end and the other end open and attached (bonded) to the face of a flange attached to the steel cylinder for pressure containment.\nThe hydraulic fluid is a matter of choice and can be either a water-antifreeze mix or hydraulic mineral oil.\nIn the method and system of the invention, an Indirect Pressurization Facility (IPF), consists of gas loading and unloading conduits, gas control mechanisms, pressure and temperature measurement devices, a bladder-equipped steel cylinder pressurization unit and a hydraulic system (pump, prime mover and surge-reservoir tank). The prime mover will be a gas engine and the pressurizer will be a centrifugal pump. Fuel for the prime mover will be natural gas from the gas well or other supply.\nStart-up operations will commence with an xe2x80x9con-offxe2x80x9d switch, when placed in the xe2x80x9conxe2x80x9d position will enable the control panel to signal the valve on the gas supply line to open and gas to flow to the suction of the indirect pressurization facility, and from there to the interior of the bladder. When the amount of gas necessary to fill the bladder is confirmed by a sensor, the fill valve will close and a signal sent to the pressure pump on the hydraulic system to pump hydraulic fluid into the annulus between the outside of the bladder and the interior of the pressurization chamber. When the annulus pressure reaches the pre-set discharge pressure, pumping into the annulus will be discontinued. At this point, the bladder will have collapsed which will compress the gas trapped inside the bladder after which the pressurized gas will be transferred to the gas sales line, or to storage. Following discharge of the gas to sales, the pressurization chamber is de-pressurized by opening the annulus to a reservoir located beneath the steel cylinder which will receive all of the fluid from the annulus at atmospheric pressure and supply the fluid to the suction of the centrifugal pump to pressurize it to the maximum working pressure of the pump and have it ready to commence the next pressurization cycle as soon as the bladder fills with gas from the supply well. The cycle is then repeated as long as the gas well has sufficient producing capacity to load the bladder at atmospheric pressure."}
{"text": "The present invention relates to a hydraulic circuit control system for a construction machine in which an operating system of the construction machine, particularly a control lever device, comprises a joystick device of the type generating an electrical operational signal (electric signal) depending on an input amount upon shift of a control lever, and a flow control valve is controlled with the operational signal for controlling the operation of an actuator.\nIn recent construction machines, particularly in those machines that are employed for various kinds of works because of convenience in use as represented by hydraulic excavators, operability has become increasingly valued in making the machines adaptable for a variety of usages. Stated otherwise, taking a hydraulic excavator as an example, the machine must be able to operate a working device as intended by an operator over a wide range from work in which primary importance is put on the amount of work carried out by the machine, e.g., excavation, to work in which fine adjustment is required in operation, e.g., leveling. To that end, it has been proposed to employ a hydraulic circuit control system in which a control lever device comprises an electric joystick for generating an electrical operational signal depending on an input amount upon shift a control lever, and the operational signal is electrically processed to control a flow control valve with a processed signal. Several known examples of such a control system are as follows.\n(1) Japanese Patent No. 2509311 entitled xe2x80x9cWorking Device Control Method for Construction Machinexe2x80x9d\nThis publication discloses a working device control method for a construction machine comprising a hydraulic control valve (operational valve), which is operated through a controller upon manipulation of an electrical lever, and a pump varying device. Modulation control is performed to absorb shocks caused upon operation of the operational valve and the pump varying device by setting a modulation pattern for rise/fall of a circuit pressure and increase/decrease of a pump delivery rate upon operation of the operational valve to restrict a maximum operating speed of the operational valve (maximum change rate of an operational signal) so that a rate of the rise/fall of the circuit pressure and increase/decrease of the pump delivery rate is gradually changed in multiple stages with a working time, and by operating the operational valve and the pump varying device so as not to move faster than the speeds set by the modulation pattern when the circuit pressure rises and falls at a constant rate with a working time. Furthermore, a cavitation is prevented from occurring upon operation of the pump varying device. This publication also discloses that a plurality of modulation patterns for the operational valve are prepared and one of the patterns is set depending on the working condition automatically or manually with selection by an operator.\n(2) JP,B 7-107279 entitled xe2x80x9cWorking Device Control Method for Construction Machinexe2x80x9d\nThis publication discloses an improvement of the modulation control in the above-mentioned (1). At the time when an electrical lever is manipulated from a shift position on the side in one direction toward the side in an opposite direction in a continuous manner and an operational signal from the electrical lever enters the opposite direction side beyond a dead zone corresponding to a neutral position, the modulation pattern having been effective so far is released and another modulation pattern for the opposite direction side is made effective. The operation of a working device and an operating feeling in the lever-reversed operation are thereby matched with each other.\n(3) JP,A 10-37247 entitled xe2x80x9cOperation Control Device and Operation Control methodxe2x80x9d\nThis publication discloses a hydraulic circuit controller for controlling the operation of a working device of a construction machine through a flow control valve, wherein a maximum change rate of an operational signal for the flow control valve is restrained to be not larger than a setting value, and the operation of the working device is controlled by changing the setting value depending on an input amount upon shift of a control lever.\nMeanwhile, there is also known a hydraulic circuit control system in which an actuator speed is controlled by controlling a delivery rate of a hydraulic pump with an operational signal instead of controlling a flow control valve with the operational signal, and a maximum operating speed of a pump displacement varying mechanism is restrained. Several examples of such a hydraulic circuit control system are as follows.\n(4) JP,B 62-13542 entitled xe2x80x9cController for Hydraulic Circuitxe2x80x9d\nThis publication discloses a hydraulic circuit controller for a closed circuit system wherein an actuator speed is controlled to a speed instructed by an operating device by controlling a delivery rate of a hydraulic pump (position of a pump displacement varying mechanism). When an operating speed of the pump displacement varying mechanism is restrained to be not larger a setting maximum speed, the setting maximum speed is changed depending on an input amount upon shift of a control lever, thereby controlling acceleration/deceleration of an actuator.\n(5) JP,B 62-39295 entitled xe2x80x9cControl System for Hydraulic Circuit Apparatusxe2x80x9d\nThis publication discloses that the controller of the above-mentioned (4) is modified so as to detect a condition of the operating device (control lever) instructing the operation to be stopped or made in the reversed direction, and to set the setting maximum speed larger than that in acceleration.\nThe above-described prior art however has the following problems.\nFirst problem: The setting value for restricting the maximum operating speed of the operational valve (flow control valve) (i.e., the maximum change rate of the operational signal) is not set corresponding to individual operating status, i.e., acceleration, deceleration/stop, and lever-reversed condition. Therefore, the operational valve cannot be always controlled at an optimum maximum change rate adapted for the operating status of a construction machine.\nSecond problem: In the lever-reversed operation, the dead zone in the vicinity of a neutral position of the flow control valve is not appropriately handled or not handled at all. When quickly reversing the control lever, therefore, the actuator undergoes a shock or stalls in the vicinity of the neutral position, causing the operator to feel a pause in the operation.\nThird problem: Since the maximum change speed of the operational valve is just restrained to the fixed modulation pattern regardless of the input amount upon shift of the control lever, an appropriate acceleration/deceleration feeling corresponding to the lever shift amount cannot be provided.\nMore specifically, in Japanese Patent No. 2509311 and JP,B 7-107279, the modulation patterns are set for the maximum operating speed of the operational valve in acceleration and deceleration/stop, and in the lever-reversed operation, the maximum operating speed of the operational valve is restricted in accordance with the modulation pattern for deceleration/stop. However, the lever reversing is performed when it is required to quickly change the moving direction of the working device in the case of, e.g., dropping mud from a bucket, bumping a boom against a vertical surface, or avoiding a risk, and a rapid response is demanded until the working device changes the moving direction. Accordingly, restricting the maximum operating speed of the operational valve in the lever-reversed operation in accordance with the modulation pattern for deceleration/stop cannot be the as providing an optimum maximum operating speed for the lever-reversed operation, and hence cannot change the moving direction of the working device with a good response (first problem).\nAlso, according to JP,B 7-107279, as soon as the operational signal indicates a reversed direction, the modulation control performed so far is ceased and another modulation control adapted for the reversed direction is started for the purpose of improving response in the lever-reversed operation disclosed in Japanese Patent No. 2509311. Taking into account a delay in the operation of the actuator responsive to the operational signal, therefore, the actuator is brought into an uncontrolled state at the moment when the operating direction is changed, which leads to a possibility that a substantial shock may occur until the moving direction of the actuator is completely changed (second problem).\nFurther, in Japanese Patent No. 2509311 and JP,B 7-107279, because the modulation pattern is fixed and the maximum operating speed of the operational valve is always restricted to the fixed modulation pattern regardless of the input amount upon shift of the control lever, an appropriate acceleration/deceleration feeling corresponding to the lever shift amount cannot be provided (third problem). In the case of returning the control lever, for example, when the control lever is manipulated so as to operate the operational valve at a speed higher than that set by the modulation pattern, the maximum operating speed of the operational valve is determined by the fixed modulation pattern regardless of a manner in which the control lever is returned, and therefore cannot be adjusted.\nIn JP,A 10-37247, since the maximum operating speed of the operational valve is not set depending on the operating status of the construction machine, the operational valve cannot be controlled at an optimum maximum change rate adapted for the operating status (first problem), and an appropriate acceleration/deceleration feeling corresponding to the lever shift amount cannot be provided (third problem). Furthermore, no consideration is paid on how to handle the lever-reversed operation (second problem).\nIn JP,B 62-13542 and JP,B 62-39295, the position of the pump displacement varying mechanism is controlled in response to an instruction from the operating device to control the pump delivery rate, thereby controlling the actuator speed. That is to say, these are not intended to control the operation of the working device of the construction machine through the flow control valve. Also, in the system of JP,B 62-39295, a plurality of maximum change rates of the operational signal are set as a function of the operational signal. However, because a control target of the control lever is the pump displacement varying mechanism, no consideration is paid to the dead zone in the vicinity of the neutral position of the flow control valve. Accordingly, if the disclosed arrangement is applied to a hydraulic circuit control system for controlling an actuator speed through a flow control valve, the maximum change rate of an operational signal is restrained in a similar manner even when the flow control valve is within the dead zone in the vicinity of its neutral position, whereby an actuator stalls for a certain period of time, causing the operator to feel a pause in the operation (second problem).\nA first object of the present invention is to provide a hydraulic circuit control system for a construction machine of the type controlling a flow control valve with an electrical operational signal to control the operation of an actuator, the control system being able to control the flow control valve at an optimum maximum change rate in any operating status of acceleration, deceleration/stop, and lever-reversed condition with resulting characteristics cited below:\n(a) in acceleration/deceleration, the machine undergoes a less shock and an operator feels no delay in the operation even with the operator manipulating a control lever quickly;\n(b) in moderate acceleration/deceleration, the actuator is moved as intended by the operator;\n(c) in stop operation, the machine undergoes a less shock and the operator feels no delay in motion toward stop even with the operator manipulating the control lever quickly; and\n(d) in quick lever reversing, the actuator can be rapidly reversed in motion.\nA second object of the present invention is to provide a hydraulic circuit control system for a construction machine, which carries out, in addition to the above, proper processing for a dead zone in the vicinity of a neutral position of the flow control valve in the lever-reversed operation, whereby the machine undergoes a less shock and the operator feels neither a delay in the operation nor a pause in the operation in the vicinity of the neutral position when the control lever is quickly reversed.\nA third object of the present invention is to provide a hydraulic circuit control system for a construction machine, which can give the operator an appropriate feeling in acceleration and deceleration corresponding to an input amount upon shift of the control lever.\n(1) To achieve the above first object, the present invention provides a hydraulic circuit control system for a construction machine comprising a hydraulic actuator for driving a working device, a hydraulic pump driven by a prime mover and producing a pressurized hydraulic fluid, a flow control valve disposed between the hydraulic actuator and the hydraulic pump and controlling a flow rate of the hydraulic fluid, and operational signal generating means for generating an electrical operational signal to instruct a flow rate of the hydraulic fluid flowing through the flow control valve, the system computing a control signal while restraining a change rate of the operational signal to be kept not more than a preset maximum change rate, and controlling the flow control valve in accordance with the computed control signal, wherein the system comprises first determining means for determining the operating status of the construction machine based on the operational signal; and first processing means for setting therein an optimum maximum change rate of the control signal for the flow control valve beforehand for each operating status of the construction machine, determining an optimum maximum change rate adapted for the operating status of the construction machine at that time based on a determination result of the first determining means, and setting the determined optimum maximum change rate as a maximum change rate of the control signal for the flow control valve.\nThus, since the first determining means determines the operating status of the construction machine and first processing means determines an optimum maximum change rate adapted for the operating status of the construction machine at that time based on a determination result of the first determining means and then sets the determined optimum maximum change rate as a maximum change rate of the control signal for the flow control valve, the change rate of the control signal for controlling the flow rate through the flow control valve is restrained to be kept not more than the determined optimum maximum change rate. Therefore, the flow control valve can be controlled at the optimum maximum change rate in any operating status of acceleration, deceleration/stop, and lever-reversed condition with such resulting characteristics as (a) in acceleration/deceleration, the machine undergoes a less shock and an operator feels no delay in the operation even with the operator manipulating a control lever quickly; (b) in moderate acceleration/deceleration, the actuator is moved as intended by the operator; (c) in operation for stop, the machine undergoes a less shock and the operator feels no delay in the motion toward stop even with the operator manipulating the control lever quickly; and (d) in quick lever reversing, the actuator can be rapidly reversed in motion, whereby working efficiency and safety are improved.\n(2) To achieve the above second object, according to the present invention, in the hydraulic circuit control system for a construction machine of the above-mentioned (1), the system further comprises second determining means for determining whether a value of the control signal for the flow control valve is within a neutral zone; and second processing means for computing the control signal in accordance with the operational signal when the value of the control signal for the flow control valve is within the neutral zone, instead of executing the processing to restrain the change rate of the control signal in accordance with the maximum change rate.\nWith those features, proper processing for a dead zone in the vicinity of the neutral position of the flow control valve is executed in the lever-reversed operation so that, when the control lever is quickly reversed, the machine undergoes a less shock and the operation can be performed without causing the operator to feel neither a delay in the operation nor a pause in the operation in the vicinity of the neutral position. As a result, operability in the lever-reversed operation is greatly improved.\n(3) In the above-mentioned (1), preferably, the first determining means determines, based on a state of the operational signal, in which one of acceleration, deceleration/stop, and lever-reversed condition the operating status of the hydraulic excavator is, and the first processing means determines the optimum maximum change rate adapted for the operating status of the construction machine at that time based on the optimum maximum change rate of the control signal set beforehand for each operating status of acceleration, deceleration/stop, or lever-reversed condition.\nWith those features, as with the above-mentioned (1), the flow control valve can be controlled at the optimum maximum change rate in any operating status of acceleration, deceleration/stop, and lever-reversed condition.\n(4) Also, in the above-mentioned (1) or (3), preferably, the first determining means determines the operating status of the construction machine based on the operational signal and a previously outputted control signal for the flow control valve.\nWith that feature, the first determining means can determine the operating status of the construction machine including acceleration, deceleration/stop, and lever-reversed condition.\n(5) To achieve the above third object, according to the present invention, in any one of the above-mentioned (1), (3) and (4), the optimum maximum change rate of the control signal for the flow control valve is set beforehand as a function of the operational signal for each operating status of the construction machine, and the first processing means computes the optimum maximum change rate based on the function of the operational signal corresponding to the operating status determined by the first determining means and the operational signal at that time.\nWith those features, the optimum maximum change rate of the control signal is set depending the value of the operational signal, and hence an appropriate feeling in acceleration and deceleration corresponding to the input amount upon shift of the control lever can be provided.\n(6) In any one of the above-mentioned (1), (3) and (4), preferably, the optimum maximum change rate of the control signal for the flow control valve is set beforehand as a function of the operational signal or a function of the previously outputted control signal for the flow control valve for each operating status of the construction machine, and the first processing means computes the optimum maximum change rate based on the function of the operational signal corresponding to the operating status determined by the first determining means or the function of the previously outputted control signal for the flow control valve and the operational signal at that time or the previously outputted control signal for the flow control valve.\nWith those features, the optimum maximum change rate of the control signal is set depending both the value of the operational signal and the previously outputted control signal, and hence an appropriate feeling in acceleration and deceleration corresponding to the input amount upon shift of the control lever can be provided."}
{"text": "This invention is directed to a fiber optic bulkhead fitting wherein the optical fiber can be optically coupled through and soldered to a bulkhead.\nOptical fibers are suitable for the transmission of digital information and are being increasingly used for that purpose. Optical fibers are superior to the employment of electrical conductors for a number of data transmission needs. However, there has been a problem in passing an optical fiber from one environment, through a bulkhead, to another environment because of the problems of reliably sealing the optical fiber to the bulkhead. Such sealing is necessary to provide sure separation between the spaces separated by the bulkhead. Accordingly, there is need for a bulkhead fitting which is compatible with both the optical fiber and the bulkhead to provide complete and reliable sealing.\nPrevious attempts at running an optical fiber through a bulkhead and sealing it with respect thereto have been unsatisfactory because of the microbending losses. When epoxy is used to seal an optical fiber with respect to a structure like a bulkhead, the forces on the optical fiber cause local changes in the index of refraction which is the microbending loss. With epoxy, the loss is about 1 decibel. This is a significant loss in such a system and cannot be tolerated in any but short, simple systems."}
{"text": "1. Field of Invention\nThe present invention relates to a method and a system for controlling a paper machine, wherein a dryer is controlled by predicting the moisture percentage of a web at a dryer part inlet and also predicting the dryer's steam pressure according to the predicted moisture percentage.\n2. Description of Prior Art\nFIG. 1 is a schematic view showing the configuration of a typical paper machine. In the figure, raw pulp is discharged from a stock inlet 81 to a wire part 82. The wire part 82 is moved in the direction of arrow A by means of rotating rolls 821. The raw pulp discharged onto the wire part 82 is subjected to drainage so as to form a web (that is paper). The web thus formed is transferred to a press part 83 for further water drainage.\nThe web subjected to water drainage at the press part 83 is transferred to a pre-dryer 84. A multitude of steam drums 841 are disposed in the pre-dryer 84 and heated by steam introduced thereinto. The web is wound around the steam drums as it is moved forward, so that the web is drived until a given moisture percentage is reached.\nThe dried web is subjected to a sizing process, such as application of a sizing agent (coating agent) at a size press 85; is further dried by an after-dryer 86; and is then take up as a product indicated by numeral 87. It should be noted that the after-dryer 86 is configured in the same way as the pre-dryer 84.\nNumerals 88 and 89 denote BM systems, both of which detect the basis weight, moisture percentage, and other data items of the web as it comes out of the pre-dryer 84 and after-dryer 86, respectively. The values of data items thus detected are input to a control apparatus not shown in the figure. The control apparatus controls the amount of raw pulp discharged onto the wire part 82 or the amount of steam introduced into the steam drums of the pre-dryer 84 and after-dryer 86, as well as the machine speed and other parameters, so that the product in question complies with predetermined specifications. Grade change control whereby different types of product are produced is also practiced commonly.\nIn grade change control, any product obtained during the time of grade change, wherein a switch is made to another type of product, will be treated as broke, i.e., non-standard paper. Therefore, the duration of grade change must be minimized in order to increase operation efficiency. To solve this problem, an invention of a method of predicting a steam pressure setpoint after grade change by simulation is described in the specification of U.S. Pat. No. 3,094,798.\nNow, the aforementioned invention is described briefly.\nThe invention described in the specification of U.S. Pat. No. 3,094,798 uses an iron model wherein the steam drums of the pre-dryer 84 and after-dryer 86 are simplified into a planar form. In the model, the state of contact among the steam drum, web, and canvas wound continuously round the steam drums is classified into five patterns. Then, the heat-transfer differential equation of each pattern is derived and converted to a difference equation, so that a steam pressure setpoint after grade change is predicted by solving the difference equation.\nFor convenience, the numbering of the equations in this specification are 5-13 and 18-23. The numbers 1-4 and 14-17 have been omitted.\nThe heat-transfer differential equations of a pattern wherein the steam drum, web and canvas are in contact with each other in this order are represented as equations 5 to 7 below.                                                         L              D                        ·                          ρ              D                        ·                          C              D                                ⁢                                    ⅆ                                                T                  1                                ⁡                                  (                  t                  )                                                                    ⅆ              t                                      =                                            h              S                        ·                          (                                                                    T                    S                                    ⁡                                      (                    t                    )                                                  -                                                      T                    1                                    ⁡                                      (                    t                    )                                                              )                                -                                    h              DW                        ·                          (                                                                    T                    1                                    ⁡                                      (                    t                    )                                                  -                                                      T                    2                                    ⁡                                      (                    t                    )                                                              )                                                          (        5        )                                                                                                                          L                    W                                    ·                                      ρ                    W                                    ·                                      C                    W                                                  ⁢                                                      ⅆ                                                                  T                        2                                            ⁡                                              (                        t                        )                                                                                                  ⅆ                    t                                                              =                            ⁢                                                                    h                    DW                                    ·                                      (                                                                                            T                          1                                                ⁡                                                  (                          t                          )                                                                    -                                                                        T                          2                                                ⁡                                                  (                          t                          )                                                                                      )                                                  -                                                                                                      ⁢                                                                    -                                          h                      WC                                                        ·                                      (                                                                                            T                          2                                                ⁡                                                  (                          t                          )                                                                    -                                                                        T                          3                                                ⁡                                                  (                          t                          )                                                                                      )                                                  ⁢                                  Evapo                  ⁡                                      (                                                                  T                        2                                            ,                                              T                        W                                                              )                                                                                                          (        6        )                                                                    L              C                        ·                          ρ              C                        ·                          C              C                                ⁢                                    ⅆ                                                T                  3                                ⁡                                  (                  t                  )                                                                    ⅆ              t                                      =                                            h              WC                        ·                          (                                                                    T                    2                                    ⁡                                      (                    t                    )                                                  -                                                      T                    3                                    ⁡                                      (                    t                    )                                                              )                                -                                    h              a                        ·                          (                                                                    T                    3                                    ⁡                                      (                    t                    )                                                  -                                                      T                    a                                    ⁡                                      (                    t                    )                                                              )                                                          (        7        )            \nThe meanings of the parameters included in equations 5 to 7 are as follows.\nLD:Drum thickness (m)LW:Web thickness (m)LC:Canvas thickness (m)Ts:Steam temperature within drum (° C.)T1:Drum's surface temperature (° C.)T2:Web (paper) temperature (° C.)T3:Canvas temperature (° C.)Ta:Dry-bulb temperature of air within hood (° C.)CD:Drum's specific heat (kJ/(kg · ° C.))CW:Web's (paper's) specific heat (kJ/(kg · ° C.))CC:Canvas' specific heat (kJ/(kg · ° C.))ρD:Drum's density (kg/m3)ρW:Web's (paper's) density (kg/m3)ρC:Canvas' density (kg/m3)hS:Coefficient of heat transfer between steam within drum and drumsurface (kJ/(m2 · sec · ° C.))hDW:Coefficient of heat transfer between drum surface and web(kJ/(m2 · sec · ° C.))hWC:Coefficient of heat transfer between web surface and canvas(kJ/(m2 · sec · ° C.))ha:Coefficient of heat transfer between canvas and air within hood(kJ/(m · sec · ° C.))\nFIG. 2 is a table that summarizes the above-listed parameters.\nThe term Evapo(T2, TW) in equation 6 is a function representing the amount of heat of evaporation removed from the web as the result of moisture evaporation, and is given by equation 8 below.\n Evapo(T2, Tw)=V(MPABS)·K·(P(T2)−P(Tw))·SB(T2)(kJ/(m2·sec))  (8)\nwhere\n    P(T)=Saturation vapor pressure (kPa) at temperature T (° C.)    SB(T)=Heat of evaporation (kJ/H2Okg) at temperature T (° C.)    TW=Wet-bulb temperature of air within hood (° C.)    V(MPABS)=Function representing moisture evaporation intensity at absolute moisture percentage MPABS, where 0.0≦V(MPABS)≦1.0 (dimensionless)    K=Drying rate coefficient (H2Okg/(m2·sec·kPa)).\nAlthough heat-transfer differential equations for patterns of contact other than those mentioned above are also given by the invention described in the specification of U.S. Pat. No. 3,094,798, these equations are omitted here to avoid complication.\nIn differential equations 5 to 7 discussed earlier, a length of time is segmented into time intervals Δt, which is determined by the machine speed, circumference of a steam drum, and other data items, so that a difference equation is derived and the numeric solution thereof is obtained. Since the web moves from the upstream side to the downstream side of the paper machine as time elapses, it is possible to calculate the web temperature at the steam drum by numerically solving the difference equation.\nFrom equation 8, EvapoMP(T2, TW)(H2Okg/(m2·sec)), which is the amount of moisture evaporated from the web per unit area and unit time, can be represented by equation 9 below.EvapoMP(T2, TW)=V(MPABS)·K·(P(T2)−P(Tw))(H2Okg/(m2·sec))  (9)\nBy using this equation, it is possible to calculate the absolute moisture percentage MPABS(j) (j=1, . . . , N) of the web after the lapse of the incremental time interval Δt as shown in equation 10 below.                                           MP            ABS                    ⁡                      (                          j              +              1                        )                          =                                            MP              ABS                        ⁡                          (              j              )                                -                                                                      10                  3                                ·                                  EvapoMP                  ⁡                                      (                                                                  T                        2                                            ,                                              T                        W                                                              )                                                  ·                Δ                            ⁢                                                           ⁢              t                        BD                                              (        10        )            where    BD=Bone-dry basis weight(g/m2)    Δt=Incremental time interval (sec)    MPABS(j) (j=1, . . . , N)=Absolute moisture percentage (%) at mesh division j\nFrom this absolute moisture percentage, it is possible to calculate the (relative) moisture percentage MP(j) (j=1, . . . , N) ( %) as shown in equation 11 below.                               MP          ⁡                      (            j            )                          =                              100            ·                                                            MP                  ABS                                ⁡                                  (                  j                  )                                                            1                +                                                      MP                    ABS                                    ⁡                                      (                    j                    )                                                                                ⁢                                           ⁢                      (            %            )                                              (        11        )            where    MP(j) (j=1, . . . , N)=Relative moisture percentage (%) at mesh division j\nFIG. 3 is a flowchart representing the algorithm of a steady-state simulation using equations 5 to 11. In the first step, the current operation status data, such as the current machine speed (m/min), basis weight setpoint (g/m2), and moisture percentage setpoint (%), are acquired. In the second step, the incremental time interval Δt for differential calculations is determined from the machine speed, drum's circumference, and other data items. In the third step, the steam temperature Ts(j) (j=1, . . . , N) within the drum is calculated from the current dryer steam pressure setpoint by using a saturation vapor pressure curve. Note that N is the number of mesh divisions.\nIn a further step, equations 5 to 11 and the difference equations derived therefrom are used to calculate the drum temperature T1(j) (j=1, . . . , N), web temperature T2(j) (j=1, . . . , N), canvas temperature T3(j) (j=1, . . . , N), and web's final moisture percentage MP(j) (j=1, . . . , N). In yet a further step, a judgment is made on convergence between the web's relative moisture percentage MP(N) and actual measured value MPMEASURE provided by a moisture sensor at a final cylinder. Convergence has been reached if the absolute value of the difference between MP(N) and MPMEASURE is smaller than the given value EP.\nIf convergence has not yet been reached, the drying rate coefficient K is corrected by ΔK to calculate the drum temperature, web temperature, canvas temperature, and web's relative moisture percentage once again. When convergence has been reached, the drying rate coefficient K, drum temperature T1(j), web temperature T2(j), canvas temperature T3(j), and web's moisture percentage MP(j) are fixed to their values at that moment, and the steady-state simulation ends.\nFor a dryer part consisting of pre-dryer and after-dryer parts, it is also acceptable to calculate the moisture percentage at an after-dryer outlet as the final moisture percentage. Alternatively, moisture percentages at the pre-dryer and after-dryer outlets may be defined as the final moisture percentages. In the latter case, a convergence calculation should be made for each of the dryer parts.\nIn the steady-state simulation heretofore discussed, the drying rate coefficient K is adjusted so that the absolute moisture percentage at the final cylinder is approximated to the actual measured value. Next, a simulation of steam pressure prediction is carried out, in order to predict the optimum steam pressure setpoint in an operation status after grade change. The simulation of steam pressure prediction is explained by referring to FIG. 4.\nIn the first step in FIG. 4, operation status data after grade change, i.e., the machine speed (m/min), basis weight setpoint (g/m2), and moisture percentage setpoint (%), are acquired. In the second step, the incremental time interval Δt for differential calculations is determined from the machine speed, drum's circumference, and other data items. In the third step, the steam temperature Ts(j) (j=1, . . . , N) within the drum is calculated from the current dryer steam pressure setpoint P (kPa) by using a saturation vapor pressure curve. Note that N is the number of mesh divisions.\nIn a further step, the value of the drying rate coefficient K determined in the steady-state simulation, as well as the value before grade change used in the steady-state simulation, for example, as the pre-dryer part inlet moisture percentage, is used to find the numerical solutions of equations 5 to 11 and their difference equations, thereby calculating the drum temperature T1(j) (j=1, . . . , N), web temperature T2(j) (j=1, . . . , N), canvas temperature T3(j) (j=1, . . . , N), and web's moisture percentage MP(j) (j=1, . . . , N) as the initial values for the difference equations.\nIn yet a further step, the value of the web's moisture percentage MP(N) at the final cylinder and the moisture percentage setpoint after grade change are compared, in order to judge convergence in the same way as in the case of the steady-state simulation. If convergence has not yet been reached, the dryer steam pressure setpoint is corrected by the given value Δt, and the drum temperature, web temperature, canvas temperature, and web's relative moisture percentage are calculated once again. When convergence has been reached, the values of these data items at that moment are fixed and the simulation of steam pressure prediction ends.\nIn such a paper machine as discussed above, controlling the process of drying a product is an important factor in order to produce products of consistent quality. Drying at the after-dryer 86 is particularly important since the drying process directly affects product quality. For this reason, it is necessary to precisely know the moisture percentage of a product at the dryer inlet.\nTraditionally, the moisture percentage of a product at the inlet of the after-dryer 86 has been calculated by using a measured value provided by the BM system 88 installed before the size press 85 and then applying, for example, equation 12 shown below. It should be noted that the absolute moisture percentage in the equation means the ratio of moisture weight to the bone-dry weight of a web which is a product.                               absMP          AFTIN                =                                                            BD                PRE                            ×                              absMP                PREEND                                      +                          CW              ·                                                100                  -                  S                                S                                                          BD            AFT                                              (        12        )            where    absMPAFTIN=Absolute moisture percentage (0.0 to 1.0) at after-dryer 86 inlet    absMPPREEND=Absolute moisture percentage (0.0 to 1.0) at pre-dryer 84 outlet (calculated by simulation)    BDPRE=Bone-dry basis weight (g/m2) at pre-dryer 84 outlet (measured with BM system)    BDAFT=Bone-dry basis weight (g/m2) at after-dryer 86 outlet (measured with BM system)    CW=Size's bone-dry coated weight (g/m2)    S=Moving average of size's (coating agent's) concentration (%).\nThe pre-dryer 84 outlet absolute moisture percentage absMPPREEND is evaluated as a solution given by simulating the steady state formed in the pre-dryer 84. However, a size with a concentration of 5 to 10% is coated at the size press 85 and therefore, the moisture percentage must be corrected by the amount of moisture produced by such coating.\nMore specifically, the first term BDPRE×absMPPREEND of the numerator on the right-hand side of equation 12 denotes a moisture weight (g/m2) per unit area at the outlet of the pre-dryer 84, whereas the second term CW·(100−S)/S denotes a moisture weight (g/m2) contained in the coated size per unit area. Since the sum of these two terms is the amount of moisture contained per unit area of a product at the inlet of the after-dryer 86, it is clear that the absolute moisture percentage is evaluated by dividing this amount by the bone-dry basis weight BDAFT measured with the BM system 89.\nIt should be noted that as the size's bone-dry coated weight CW, equation 12 uses the value calculated by equation 13 below, which is the difference between the bone-dry basis weights measured with the BM systems 88 and 89.CW=BDAFT−BDPRE  (13)"}
{"text": "The clutch band of the present invention relates to the class of band clutch constructions and configurations as disclosed in U.S. Pat. No. 2,518,453, issued Aug. 15, 1950 to J. M. Dodwell and particularly in U.S. Pat. No. 3,731,773, issued May 8, 1973 to M. Austin and N. Boulton. The above identified patents disclose a band-type, free wheeling clutch embodying drive and driven clutch members, one of the members being provided with a V-groove into which the narrowed end of at least one clutch band is positioned. The wide end of at least one clutch band is anchored to the other of the clutch members such that upon rotation of the drive member in one direction, the driven member is driven in the same direction through the engagement of the side edges of the narrow end of at least one clutch band with the side walls of the V-groove. When the driven member rotates at a faster speed than the drive member, at least one clutch band permits the driven member to overrun the drive member.\nIn the development of band-type, free-wheeling clutch assemblies of the above type, one of the major problems was in providing an adequate anchor at one end of the clutch band. Various structural arrangements have been proposed wherein the clutch band is formed as an elongated structure having a wider, heavier end portion to provide an adequate anchor to engage one of the clutch members and a narrower, lighter end portion extending into the V-groove of the other clutch member. In the operation of the overrunning type clutch assemblies containing at least one band, during the overrunning condition, the wider portion, being heavier, flies radially outwardly with more force than the narrow portion, being lighter, such that the wider portion pulls the narrower portion into the V-groove of the inner member thus prematurely wearing out the narrower portion of the band. Such premature wear has severely reduced the number of clutching cycles obtainable with such clutch assemblies."}
{"text": "Magnetic random access memories (MRAMs) in which the magnetoresistive effect of a ferromagnetic material is used have recently attracted interest as next-generation solid-state nonvolatile memories capable of speeding up read/write operations, increasing capacity, and realizing low-power operation. In particular, magnetoresistive elements including a ferromagnetic tunnel junction have received a great deal of attention since the finding of the ability of exhibiting a high rate of magnetoresistive change. The ferromagnetic tunnel junction has a three-layer structure including a storage layer having a variable magnetization direction, an insulating layer, and a fixed layer that faces the storage layer and maintains a predetermined magnetization direction.\nA magnetoresistive element having the ferromagnetic tunnel junction is also called a magnetic tunnel junction (MTJ) element. A write method (spin transfer torque writing) using spin-momentum transfer (SMT) has been proposed for this magnetoresistive element."}
{"text": "1. Field of the Invention\nThis invention relates to the oxidation of organic compounds. In one aspect, the invention relates to the oxidation of compounds having at least one activated methylene radical while in another aspect, the invention relates to oxidizing such a compound in a multiphase system employing a catalyst comprising a synergistic combination of elemental carbon and a phase-transfer catalyst.\n2. Description of the Prior Art\nHawthorne et al., \"Base-Catalyzed Autoxidation of 9,10-dihydroanthracene and Related Compounds\", Oxidation of Organic Compounds--I, Advances in Chemistry Series, 75, 14 (ACS 1968) teach the oxidation of 9,10-dihydroanthracene to anthraquinone by contacting the dihydroanthracene with oxygen and benzyltrimethylammonium hydroxide. Pyridine was used as a solvent.\nKang Yang, JOC, 42, 3754 (1977), teaches the oxidation of fluorene to fluorenone in the presence of charcoal and a base. The base is potassium t-butoxide in either a t-butyl alcohol or sodium hydroxide solution.\nAlveri et al., \"Autoxidation of Diarylene Methanes and Related Compounds in the Presence of Phase-Transfer Catalysts\", Tetrahedron Letters, 24, 2117 (1977), teach the use of a phase-transfer catalyst, e.g., dicetyl diethylammonium chloride, in the autoxidation of diarylene methanes. Fluorene was converted to fluorenone in 100 percent yield after 24 hours at 30.degree. C. using the catalyst, oxygen and an aqueous sodium hydroxide/benzene biphasic reaction medium.\nWhile all of the above teachings demonstrate utility, each suffers from one or more disadvantages. Principal among these disadvantages are undesirably slow reaction rates or the requirement of a neutralization step, the latter consuming starting reactants and generating a brine waste stream. Other disadvantages include the use of undesirable solvents, such as benzene, and the absence of catalyst and caustic recycle."}
{"text": "The present invention relates to an apparatus and method for the solidification of sludges by kneading sludges in soft ground layers with a hardener and solidifying the sludges.\nAs the conventional method for solidifying sludges deposited and accumulated in bottom portions of harbors, bays, rivers and lakes, there can be mentioned a method in which a hardener is added to sludge in situ and the hardener-incorporated sludge is kneaded. The reason why the sludge is solidified in situ in the deposited and accumulated state is that the amount of the water contained in the sludge is held to a minimum and the solidification treatment can be performed conveniently. If the sludge is dug out and placed on the land for solidification, the water content in the sludge is increased greatly compared with the sludge deposited naturally.\nThis known method, however, is defective in that when the deposited sludge is kneaded and agitated by a kneading machine or the sludge is dredged, sea water or the like is contaminated in a broad region to cause secondary pollutions such as the generation of bad odors. Moreover, when the sludge is agitated, water or untreated sludge flows from neighboring sludge layers into the sludge being treated, and therefore, a large quantity of a hardener must be added. In this case, the hardener supplied in such a large amount readily flows into neighboring sludge layers and is wasted. Agitation of the sludge or the like is performed by agitating blades of the kneading machine. Accordingly, the sludge being treated is not completely separated from the sludge present on or around the outer periphery of the rotation locus of the agitating blades and therefore, uniform kneading of the sludge and hardener is not attained, and it is difficult to perform the solidification treatment in good kneading conditions."}
{"text": "The present invention relates to a structure of the insulator--semiconductor type in which the semiconductor is a III-V compound and to processes for the production of said structure. It can be used on the one hand in microelectronics where it permits the production of high performance (speed, degree of integration, etc) components of the MIS type (metal--insulator--semiconductor) and on the other hand in optoelectronics where it permits the passivation of the surfaces of optoelectronic components, e.g. semiconducting lasers.\nFor some years now, an increasing importance and interest has been attached to composite semiconductors of the III-V type (e.g. GaAs or InP). These materials have in fact remarkable properties making them particularly suitable for the construction of integrated electronic components of the MIS type, where they lead to a considerable improvement in performance figures compared with conventional components based on silicon and germanium. Among these properties, particular reference is made to the high electronic mobility (approximately 10,000 cm/s for the GaAs and 5,000 for InP, as compared with only 1,200 for Si and 3,600 for Ge), the forbidden band width (approximately 1.4 eV compared with 1.1 eV for Si and 0.7 eV for Ge) and the high resistivity (approximately 10.sup.6 .OMEGA.cm).\nHowever, hitherto the hopes set on III-V compounds have not been crowned with success due more particularly to the fact that it has proved very difficult to make a suitable interface between the insulator and the semiconductor. Attempts have been made to deposit insulators such as SiO.sub.2, Al.sub.2 O.sub.3, Si.sub.3 N.sub.4, Ge.sub.2 N, GaN, etc . . . , but the electrical characteristics (current--voltage, capacitance--voltage, interface charges) obtained for these insulator--semiconductor interfaces do not provide the necessary high quality for the construction of MIS components."}
{"text": "1. Field of the Invention\nThe present invention relates to a technology for realizing low damage sputtering regardless of materials (such as, an inorganic or organic material) in a surface analysis method, and relates to a technology for realizing improvement of sensitivity by improving secondary ion yield in a secondary ion mass spectroscopy method.\n2. Description of the Related Art\nAn ion source for surface analysis that can perform sputtering without any damage to a target sample has not yet been developed. In the surface analysis, argon ion (Ar+) is the most common ion species for sputtering, but it is known that the occurrence of damage due to the sputtering is high.\nIn addition, in a secondary ion mass spectroscopy method (SIMS) as one of surface analysis methods, a primary ion beam that has been used so far is a noble gas ion or a metal ion (Cs+, Ar+, Ga+, Au+, or the like). Some of them can be reduced to a small beam in the order of several tens nanometer, but the occurrence of large damage to a sample is a common drawback.\nIn addition, if these ions are used as a primary ion source, secondary ion yield is very low, and secondary ion generation efficiency is low. Therefore, in order to overcome the drawback of the SIMS using them as the primary beam, a cluster ion SIMS has been developed. A beam source thereof is Au3++, Bi3++, or the like. By using the cluster ion (Au3++, Bi3++, or the like) consisting several atoms, desorption efficiency of the secondary ions is significantly increased in a non-linear manner. Such result is due to the generation of ablation.\nOn the other hand, because a target sample surface and its vicinity are significantly damaged, application of the conventional system to a biological material is difficult; and nondestructive observation of molecule ions is difficult; specifically, the sample receives large fragmentation, and a surface of the sample is decomposed and polymerized.\nA cluster ion source of C60+ ion is commercialized; and hence, a low damage sputtering technology is realized though in a limited manner. Further, the desorption efficiency is further increased in the SIMS using the C60+ ion source as the primary ion source. However, the following phenomena are caused: (1) an inorganic material is contaminated with a carbon component derived from C60; (2) craters are generated in a surface of the material so that surface destruction occurs; (3) a biological sample or the like is significantly damaged; and (4) the secondary ion yield is low in the SIMS, and when the beam diameter is decreased, ionic strength is weakened so that utility value as the SIMS is deteriorated (particularly in an organic material). Refer to Japanese Patent Application Laid-open No. 2005-134170, Journal of Physical Chemistry B, 108, pp 7831-7838, and Applied Surface Science 231-232, pp 936-939, FIG. 4.\nThere is a surface analysis method utilizing a gas cluster ion beam (GCIB) that has been recently popular, in which noble gas (such as argon (Ar)) is ejected in vacuum to form a jet stream, gas temperature is decreased, and neutral clusters having an n value of Arn+ of a few thousands to a few tens of thousands are formed and ionized to generate Arn+, which is accelerated to impact the sample.\nWith this method, depth profile analysis with low-damage sputtering for an organic material (such as a polymer) is confirmed to be effective and is commercialized. However, for an inorganic material (such as a ceramic material) that is relatively hard, the sputtering speed is extremely slow so that it is not practical. Therefore, a range of the sample types to be analyzed is inevitably limited to mainly organic industrial materials.\nIn addition, when the GCIB is used as the primary ion source in the secondary ion mass spectroscopy method, it is known that the secondary ion yield thereof is low; and hence, it is not practical when used for improving sensitivity in the secondary ion mass spectroscopy method. Refer to Japanese Patent Application Laid-open No. Hei 04-354865, Japanese Patent Application Laid-open No. 2008-116363, and Analytical Chemistry, 2011, 83(10), pp 3793-3800, FIG. 7.\nIn addition, an ion beam technology using a charged droplet method has been developed. In this method, a capillary is disposed in the atmosphere, solvent is supplied through inside of the capillary, and an extraction electrode that is applied with a high voltage negative with respect to the capillary is disposed in front of the capillary so as to generate ions in the atmosphere.\nA vacuum chamber is separated into several steps from low vacuum side to high vacuum side with small diameter orifices. The ions are made to pass through the orifices and are transported to vacuum atmosphere so as to be used as ion beam. In this case, the cluster ions generated in the atmosphere inevitably collide with gas molecules in the atmosphere so that many ions are scattered. Therefore, the amount of ions that are actually transported to the vacuum side and can be effectively used is small; and in addition, downsizing of the cluster ion (fission of the cluster) also occurs due to vaporization in the atmosphere side.\nIn addition, to use the ion beam, it is necessary to apply a high voltage, which is positive with respect to the ground potential, to the capillary as a source, and it is also necessary to apply a high voltage to parts for lens effect or the like in a low vacuum region during the ion transportation process. Therefore, discharge phenomenon tends to occur in various parts. Consequently, it becomes difficult to stably obtain the ion beam, and it is also difficult to decrease the beam size to be small.\nOn the other hand, a differential pumping system for evacuating the separated vacuum chamber also becomes large in scale which causes difficulty when in use. Refer to Japanese Patent Application Laid-open No. 2011-141199.\nConsequently, a practical ion source that can support various types in etching layer-by-layer without damaging a surface of the sample after irradiation has not been developed yet, and an ion source succeeding in dramatic improvement of sensitivity in the secondary ion mass spectroscopy method has also not yet been developed.\nA charged droplet ion source of the related art is described below. In FIG. 5, a charged droplet ion source 701 includes a vacuum chamber 710.\nThe vacuum chamber 710 is connected to first and second vacuum evacuating devices 729a and 729b so that the inside of the vacuum chamber 710 can be evacuated.\nAn extracting electrode 721 is provided with a small hole (orifice) so that gas flows in the vacuum chamber 710 through the extracting electrode 721 when the inside of the vacuum chamber 710 is evacuated. First, the inside of the vacuum chamber 710 is evacuated by the first and second vacuum evacuating devices 729a and 729b. \nAn emission tube (capillary) 703 is disposed outside the vacuum chamber 710.\nThe distal end of the emission tube 703 is directed towards the small hole of the extracting electrode 721; and a base part thereof on the opposite side is connected to a liquid supply pipe 743. The liquid supply pipe 743 is connected to an ionization liquid supply device 705.\nThe ionization liquid supply device 705 includes a liquid storing portion 732 and a liquid feeding pump 731. The ionization liquid stored in the liquid storing portion 732 is supplied to the base part of the emission tube 703 through the liquid supply pipe 743 by the liquid feeding pump 731, passes a thin tube in the emission tube 703, and is emitted to the outside of the emission tube 703 from an emission opening 735 at the distal end of the emission tube 703. The emission tube 703 is surrounded by an outer cylinder 707. When carrier gas (here, nitrogen gas) is supplied from a carrier gas source 708 to the inside of the outer cylinder 707, the gas is released from a distal end opening 736 of the outer cylinder 707.\nThe emission opening 735 is disposed between the distal end opening 736 of the outer cylinder and the small hole of the extracting electrode 721. Around the emission opening 735, there is formed a flow of the carrier gas from an upstream side as the base side of the emission tube 703 to a downstream side on which the extracting electrode 721 is located with the small hole.\nAn extraction power supply 728 is disposed outside the vacuum chamber 710.\nIn a state where the carrier gas supplied from the carrier gas source 708 is released from the distal end opening 736, the liquid feeding pump 731 supplies the ionization liquid to the emission opening 735, the extraction power supply 728 applies a voltage between the emission tube 703 (made of a metal here) and the extracting electrode 721 so that an electric field thereof extracts droplet cluster ions charged with a positive charge from the ionization liquid positioned in the emission opening 735. Then, the cluster ions pass through the small hole of the extracting electrode 721 and enter the inside of the vacuum chamber 710.\nOn the downstream side of the extracting electrode 721, there are disposed accelerating electrodes 722 and 723 with small holes and transport lens electrodes 724 and 725. When voltages are applied to the electrodes 722 to 725, the droplet cluster ions entering the inside of the vacuum chamber 710 pass through holes formed in the electrodes 722 to 725 so as to be a droplet cluster ion beam, and further propagates toward the downstream side.\nA size of an initial droplet cluster ion generated in the atmosphere is approximately 100 nm in diameter. However, the droplet cluster ion generated in the atmosphere is downsized due to Rayleigh fission that occurs when Coulomb repulsion of itself exceeds surface tension of the droplet. Further, the droplet cluster ions inevitably collide with gas molecules in the atmosphere so that many ions are scattered. Therefore, only a small amount of the droplet cluster ions can enter the inside of the vacuum chamber 710, and the size of the droplet cluster ion is decreased to be smaller than that of initially generated one.\nIn addition, for use as the droplet cluster ion beam, it is necessary to apply a positive high voltage with respect to the ground potential to the emission tube 703 as the generation source. Further, it is also necessary to apply high voltages to the extracting electrode 721, the first accelerating electrode 722, and the transport lens electrode 724 disposed in the low vacuum environments in the vacuum chamber 710. Therefore, an arcing phenomenon is apt to occur in the vacuum chamber 710, and hence it is difficult to obtain the droplet cluster ion beam.\nIn addition, it is necessary to separate the atmosphere outside the vacuum chamber 710 from the inside space of the vacuum chamber 710, both of which are connected to each other through the small hole of the extracting electrode 721. Therefore, the first and second vacuum evacuating devices 729a and 729b for evacuating the inside space of the vacuum chamber 710 are required to be large ones; and hence, difficulty arises when they are used in that they occupy large areas and in terms of cost.\nConsequently, in the ion source on the conventional technology, disposing the emission opening of the emission tube in the atmosphere so that the droplet cluster ion beam is generated in the atmosphere provides small amount of the droplet cluster ions that can be actually used. Hence, the conventional technology is of little practical use."}
{"text": "Hydraulic pumps compress and move hydraulic fluids by mechanical action in order to generate and transmit power. In hydraulic tool systems, hydraulic pumps provide high-pressure fluid to various actuators that transmit forces necessary to perform work. The actuators of such hydraulic tool systems often require that hydraulic fluid be provided at different flow rates and pressures for proper function. Variable displacement pumps can be utilized to accommodate the varied flow rate and pressure requirements, both individually and collectively, of the multiple actuators of hydraulic tool systems.\nSwashplate-type axial piston pumps are variable displacement pumps commonly used in hydraulic tool systems. Swashplate-type axial piston pumps include a plurality of plungers having one end held against an engagement surface of a tiltable swashplate. A ball-and-socket slipper joint is provided at the interface between each plunger end and the engagement surface of the swashplate to allow for relative sliding and pivoting motion. Each plunger reciprocates within an associated cylinder as the plungers rotate relative to the tilted engagement surface of the swashplate. When a plunger is retracted from an associated cylinder, low-pressure fluid is drawn into that chamber. When the plunger is forced back into the cylinder by the engagement surface of the swashplate, the plunger pushes fluid from the cylinder at an elevated pressure.\nEach cylinder and associated plunger together at least partially form a pumping chamber configured to intake hydraulic fluid from an inlet passage and to discharge hydraulic fluid into an outlet passage. Each pumping chamber interfaces with the inlet passage and the outlet passage through a port plate. The port plate includes an inlet port through which hydraulic fluid is drawn from the inlet passage into the pumping chamber and an outlet port through which hydraulic fluid is expelled from the pumping chamber into the outlet passage. As a plunger of a pumping chamber moves from a top-dead-center (TDC) position at the end of a discharge stroke to a bottom-dead-center position at the end of an intake stroke, the plunger passes the inlet port as it rotates relative to the port plate. As a plunger of a pumping chamber moves from a bottom-dead-center (BDC) position at the end of an intake stroke to a top-dead-center position at the end of a discharge stroke, the plunger passes the outlet port as it rotates relative to the port plate.\nThe tilt angle of the swashplate is directly related to an amount of fluid pushed from each cylinder during a single relative rotation between the plungers and the swashplate. Similarly, based on a restriction of a fluid circuit connected to the pump, the amount of fluid pushed from the cylinder during each rotation is directly related to the flow rate and pressure of fluid exiting the pump. Accordingly, a higher swashplate tilt angle of a pump equates to a greater flow rate and/or pressure of the pump, while a lower swashplate tilt angle results in a lower flow rate and/or pressure. Likewise, a higher swashplate tilt angle requires more power to produce the higher flow rates and pressures than does a lower swashplate tilt angle. As such, when the demand for fluid from the hydraulic tool system is low, the swashplate angle is typically reduced to lower the power consumption of the pump."}
{"text": "(1) Field of the Invention\nThe invention relates to an inflatable packing device including a sophisticated elastomeric inflatable bladder, either alone or in combination with a sophisticated cover.\n(2) Definition of Terms\nAs used herein and in the claims, the phrase \"inflation initiation\" refers to the location or point on the exterior of the device where first flexing of the contour of the device resulting from effective inflation is expected to occur. Inflation initiation can occur at a plurality of locations or points, depending upon choice of design.\nAs used herein and in the claims, the phrase \"inflation element\" means: the sub-assembly generally composed of the bladder, ribs, cover, upper securing means and lower collars or securing means.\nAs used herein and in the claims, the phrase \"point of contact\" means: the initial and subsequently latest expected location of interface between the exterior of the device and the wall of the well during effective inflation.\nAs used herein and in the claims, the phrase \"effective inflation\" means: the quantum of expansion of the bladder during the setting of the packing device from the run-in position of the apparatus to from between no more than about 70% to no more than about 85%, by volume, of the interior of the bladder when fully set in the well bore.\nAs used herein and in the claims, the phrase \"departure angle\" means: the angle between a straight line parallel to the longitudinal axis of the well and along the inside diameter wall of the well passing through a point of contact and a straight line drawn tangent to the exterior surface of the device for an interval of length extending from the point of contact to a distance of about one run-in diameter, this line too passing through the same point of contact. The longitudinal axis of the borehole and the two lines defining the departure angle must all be coplanar.\nAs used herein and in the claims, the phrase \"expansion profiles\" means: the transitional forms taken by the flexible portion of the inflation element during effective inflation.\nAs used herein, the phrase \"uniform inflation profiles\" means: the circumstance when the \"expansion profiles\" taken by the inflation element closely approximate straight line profiles from the point of contact to the end of the collar.\nAs used herein, the phrase \"expansion ratio\" means: the ratio of the diameter of the fully set inflation element, divided by the run-in diameter of the inflation element.\n(3) Description of the Prior Art\nInflatable packers, bridge plugs, and the like, have long been utilized in subterranean wells. Such inflatable tools normally comprise an inflatable elastomeric bladder element concentrically disposed around a central body portion, such as a tube or mandrel. A sheath of reinforcing slats or ribs is typically provided exteriorally around the bladder with an elastomeric packing cover concentrically disposed around at least a portion of the sheath. Generally, a medial portion of the sheath will be exposed and without a cover for providing anchoring engagement of the packer to the wall of the well. Pressured fluid is communicated from the top of the well or interior of the well bore by means of a down hole pump to the interior of the body and thence through radial passages provided for such purpose or otherwise around the exterior of the body to the interior of the bladder during inflation.\nNormally, an upper securing means engages the upper end of the inflatable elastomeric bladder and the reinforcing sheath (if included in the design), sealably securing the upper end of the bladder relative to the body, while a lower collar or securing means engages the lower end of the bladder and reinforcing sheath, securing the lower end of the bladder for slidable and sealable movement relative to the exterior of the body, in response to inflation forces. The elastomeric cover may be secured to the exterior of the sheath or placed around the exterior of the bladder, in known fashion.\nWith inflatable packers of this type, it has been observed that the portion of the bladder adjacent the exposed sheath section of the packer prematurely inflates prior to the other portions of the bladder which are reinforced against expansion by the reinforcing sheath and/or the elastomeric packing cover element. When the inflation element expands, one end of the bladder moves toward the other end of the device, and the bladder area adjacent the exposed sheath inflates until it meets the wall of the well bore, which may be cased or uncased. If the well bore is uncased, the well bore will have an earthen wall, and if the well bore is cased, the wall of the well bore will be the internal diameter surface of the casing.\nIt has been noted in a number of prior art designs that when service conditions encompass moderate expansion ratios, a propensity for the bladder to pinch around the exterior of the body arises, creating either a seal or a convoluted fold in the bladder that sometimes prevents the effective communication of further fluid throughout the bladder and preventing contiguous inflation propagation. The pinching seal and/or fold(s) can become entrenched in the bladder whereupon they obstruct further passage of fluid employed for inflating the bladder and therein keep fluid from reaching the farthest portions of bladder to be inflated. When this occurs in service, it always results in a soft set condition and in the imminent loss of seal between the cover and wellbore. This problem is discussed in detail in Eslinger, et al. \"Design and Testing of a High-Performance Inflatable Packer,\" SPE 37483, Society of Petroleum Engineers (1997). Tools designed to control inflation shape problems are discussed in the Eslinger paper are described in detail in U.S. Pat. No. 5,605,195 issued Feb. 25, 1997, and entitled \"Inflation Shape Control System For Inflatable Packers,\" and in U.S. Pat. No. 5,507,341 issued Apr. 16, 1996, and entitled \"Inflatable Packer With Bladder Shape Control.\"\nFolds in the bladder can be expected to occur in prior art devices like that shown in FIG. 18 when the expansion ratio is greater than 2:1. Designs of this sort inherently experience large departure angles and unfavorable expansion profiles when the expansion ratio is about 2:1 or more. By utilization of the design of the present invention, the departure angle is preferably controlled at no more than about 15.degree. and the inflation element experiences a uniform inflation profile and therefore, no folds or pinches will occur even if the expansion ratio is 3:1, or even higher. Elimination of the propensity to form folds and pinches in the present invention can be attributed to exceptionally low departure angles throughout inflation and the propagation of uniform inflation profiles throughout effective inflation.\nThe formation of folds creates unusually high triaxial stresses and strains in the vicinity of the fold. Correspondingly, these triaxial stresses and strains create a condition that causes localized failure of the bladder by means of cracking and/or tearing. Failure occurs because the physical properties of the elastomeric material composing the bladder are not adequate to survive the localized triaxial stresses and strains.\nExcept for the devices described in my patents U.S. Pat. No. 5,469,919, U.S. Pat. No. 5,564,504 and U.S. Pat. No. 5,813,459, all other prior art devices having an element construction similar to that shown in FIG. 18 experience large departure angles and unfavorable expansion profiles when the expansion ratio is greater that 2.00:1, i.e., departure angles greater than 25.degree. at a 2:1 expansion ratio and expansion profiles similar to that shown in FIG. 18. An expansion profile would be deemed unfavorable if the slope of the exterior surface at any point on the inflation element exceeds 15.degree. relative to the longitudinal axis of the wellbore. The term \"unfavorable expansion profile\" is only applicable to the \"effective inflation\" portion of the inflation cycle. The propensity to form pinching seals and folds is directly related to undesirable combinations of expansion ratio, departure angles and expansion profiles of the device. In prior art devices, pinching seals and folds are experienced upon the combination of departure angles greater than about 15.degree. and an expansion ratio greater than about 2.35:1.\nWith regard to covers, at expansion ratios of 2:1 and more, the departure angle in prior art devices other than those for the preferred embodiments in my aforementioned patents will be greater than 20.degree. and the combination of a departure angle greater than 20.degree. and an expansion ratio greater than about 2:1 has been observed to result in cracking and tearing in covers. Once a tear or tears occur, non-uniform rib spacing results. Non-uniform load distribution within the cover also occurs and general discontinuity of the cover results. These conditions, in turn, can result in extrusion of the bladder between ribs resulting in subsequent failure of the bladder and service failure of the device.\nIn my U.S. Pat. Nos. 5,469,919, and 5,564,504, and 5,813,459 entitled \"Programmed Shaped Inflatable Packer Device,\" issued Sep. 29, 1998, I disclose methods to abate the formation of pinching seals and folds during inflation of prior art devices by using a design which includes a series of shaped-controlling means on an elastomeric packing cover along the length of the bladder in the form of high and low modulus modules of varying lengths and thicknesses. While this design is an advancement in the art, the design of the modules leaves comparatively sharp angled transitional chamfers and significant size Differences between the high and low modules. These chamfers and different diameters are of such magnitude that they are easily detected by the naked eye. The short transitional chamfers give rise to localized stresses and strains in expanded covers. These localized stresses and strains can cause cracking and/or tearing in the covers which can ultimately result in device failure.\nIn another prior art device which was subjected to service conditions having expansion ratios of 2.35:1 and 3:1, the minimum achievable departure angles were about 15.degree. and 23.degree., respectively. This device used a plateau cover interval concept in accordance with my patents U.S. Pat. No. 5,469,919, U.S. Pat. No. 5,564,504 and U.S. Pat. No. 5,813,459 and has been made commercially available by High Pressure Integrity, Inc. under the product name \"Z-44\". While this product was an advancement and improvement over other prior art devices, the variations of constant thickness cover intervals with abrupt and relatively short transitions from one thickness to another caused comparatively high localized stress and strain concentrators in the cover which occasionally resulted in cracking and tearing of the cover. Z-44 and similar devices always exhibited rib kinking and experienced occasional rib cutting of the bladder. Additionally, inflation profiles exhibited plateau intervals (intervals of constant diameter along the length of the device) rather than relatively straight sloped profiles in the interval between the last point of contact with the casing (POC) and the end of the collar. Additionally, the plateau cover interval concept abated the formation of pinches and folds in bladders at moderate expansion ratios, but did not eliminate their occurrence at expansion ratios greater than 2.35:1.\nThe ability to successfully deflate and retrieve an inflatable device is a common service requirement. A pinch or fold might still have formed in a bladder during inflation even though the inflation element effected a satisfactory seal against the wall of the well. During deflation, a fold can pinch and seal around the body, obstructing the transmission of fluid out of the lower portions of the bladder and thereby prevent complete deflation of the bladder. Once a fold is formed, it is permanently entrenched in the bladder and results in multiple layers of bladder beneath the ribs. These layers in turn result in a deflated diameter which is greater than the initial run-in diameter of the inflation element. Retrieval of the device to the earth's surface is thus compromised since the device might not be able to pass through restrictions in the well bore as it is moved upwardly therein.\nI have now discovered that the problems described above can be further abated by providing an inflatable packing device having a combination of an excellent uniform expansion profile during effective inflation and minimal departure angles throughout the inflation cycle.\nThe invention permits orchestration of varying sophisticated contours and configurations in the bladder or in a combination of the bladder and the cover to provide a uniform expansion profile in an expected, i.e., pre-determinable, manner which can be achieved with only minimal or nominal experimentation which will be within the ordinary skills of those knowledgeable in the design and use of inflatable elastomeric devices for use in subterranean wells, and by adhering to the teachings herein."}
{"text": "This invention relates to an improved method of preparing high purity polyamic acids, to be transformed into polyimide resins which, in turn, are used in the processing of circuitry in the microelectronics industry. The improvement relates to increased yield.\nDue in part to such characteristics as high tensile properties, desirable electric properties and excellent stability in the face of heat and water, polyimide resins have been widely used in the processing of microelectronic circuitry. However, the electrical properties of microelectronic circuits can be impaired by the presence of ionic impurities, hence it is vital that polyamic acids and polyimides prepared from them, be ionically \"pure\". The necessity of using polyamic acids with reduced ionic impurities, when working in the electronics field, is well-recognized. For example, see U.S. Pat. No. 4,225,702, issued Sept. 30, 1980, to Mikino et al.\nMikino et al. disclose a method of preparing a polyamic acid with reduced ionic impurities. It involves reacting a diamine or diaminoamide compound monomer with a polycarboxylic acid dianhydride monomer, wherein said monomer compounds have been purified by a recrystallization process so that they have reduced ionic impurities.\nWhile recrystallization is an acceptable method of reducing impurities, it has the disadvantage of yielding purified monomers in an amount equivalent to only about 80 to 85% by weight of the unpurified monomer.\nThe present invention provides an improved method of preparing high purity polyamic acids that provides an increased yield of the purified monomers."}
{"text": "This invention is directed to disposable swimpants and swimsuits for incontinent adults and children. More particularly, the swimpants include a material which is permeable to fluid, but substantially impermeable to larger bowel movement material.\nFor example, disposable swimpants and swimsuits for pre-toilet trained children usually have absorbent cores and moisture barriers to prevent leaks of urine and bowel movements. When a child swims while wearing a disposable swimpant, water gets inside the swimpant. One potential problem is that if bowel movement material is also inside the swimpant, when the child leaves the water, or stands up above the surface of the water, water will exit the swimpant through the leg openings and/or the waist opening and the bowel movement material may also exit along with the water, thus creating a sanitation problem. Even if the disposable swimpant has containment flaps, the bowel movement material could conceivably exit the swimpant along with the water through the leg openings and/or waist opening.\nAnother possible problem is that sand may end up within the swimpants and swimsuits of children at the beach. The sand can be an irritant to the skin and uncomfortable to the wearer. However, currently available disposable swimwear is not designed to allow sand to move away from direct contact with the wearer's skin.\nThere is a need or desire for a swimwear garment that provides bowel movement containment before and after swimming and also allows sand to move away from direct contact with the wearer's skin."}
{"text": "1. Technical Field\nEmbodiments of the present invention relate to a projector and a keystone distortion correction method for a projector.\n2. Related Art\nIn a projector, the tilt of the chassis in the installation condition causes a keystone distortion in the projected image. Although such a keystone distortion has been corrected by the user with button operations, automated correction methods have gradually been adopted in recent years.\nFor example, there has been known a method of installing a range sensor for measuring the distance between the projector and the screen in the projector, detecting the tilt angle of the projector with respect to the screen based on an output of the range sensor, and correcting the keystone distortion in accordance with the tilt angle (see, e.g., JP-A-2000-122617).\nFurther, there has also been known a method of providing an acceleration sensor inside the projector to detect the tilt of the projector, and thus performing the keystone distortion correction of the projection image (see, e.g., JP-A-2003-283963).\nHowever, according to these methods, an expensive component such as a range sensor or an acceleration sensor needs to be incorporated therein, and a problem of increase in cost arises. Therefore, particularly in projectors, it is difficult for low price popularization models to incorporate it from a viewpoint of cost, and accordingly, a problem arises that the automated keystone distortion correction function (auto keystone function) can hardly be implemented in popularization models."}
{"text": "Siphons supplying carbonated water (soda) use dispenser valve heads which, normally, comprise a valve body crossed by a circulation channel. This channel, at one end, ends in a dispenser outer spout, while at the other ends into a dip tube located inside the bottle.\nIn general, the valve body is formed by the head walls, therefore this kind of head has a large volume and require a large mass of material for its manufacture.\nFurther, in order to allow dispensing of the liquid contained into the bottle, a check valve is included into this circulation channel, which is actuated by means of a rigid lever projecting outwardly from the valve head. Therefore, this kind of rigid lever is the means used typically by conventional valve heads.\nOn the other hand, by means of the improvements of the instant invention, a head is obtained the valve body of which is covered by a covering hood having a diametrical slot, constituting an operating channel for command thereof. The latter, being intended to hold a check valve, has its fulcrum in said covering hood and comprises an articulated arm which is foldable over the covering hood, wherein folding stops and a folding straps are located for retaining the end of the arm which partially closes said diametrical slot.\nThis constitution of the head affords a better operation and requires a smaller mass of constituting material, which makes the head lighter and less expensive. In fact, the valve body does not act as a support for the command means, therefore, the latter is reduced to its minimum possible expression. This simplification of the valve body is compensated by the presence of a covering hood protecting it from shocks and dirt.\nFurther, the instant covering hood is designed for serving as mounting means for a command means the articulated arm of which is foldable on said hood, where it has folding stops and a strap retaining its end. The inclusion of a diametrical slot allows fixing the strap and unfolding the articulated arm.\nThe smaller base of the hood and the cited stops prevent the folding articulated arm from being pressed in by any accidental pressure, in which case the valve stem would be displaced and an undesirable delivery of liquid would take place. Therefore, frequent shocks and pressures produced during storage and loading and unloading of trucks may not actuate the command means which is folded inside the diametrical slot and also properly strapped.\nThe perfect compatibility between the forming elements is demonstrated by the engagement between the valve body and the covering hood, since they not only have connecting means therebetween, but, further, the hood has a side mouth. This mouth, originated at the edge of the open larger base, allows the passage of the outer spout through which the carbonated liquid is delivered."}
{"text": "The present invention concerns a model of the acoustic sound channel associated with the human phonation system and/or music instruments and which has been realized by means of an electrical filter system.\nFurthermore, the invention concerns new types of applications of models according to the invention, and a speech synthesizer applying models according to the invention.\nThe invention also concerns a filter circuit for the modelling of an acoustic sound channel.\nIn its most typical form, this invention is associated with speech synthesis and with the artificial producing of speech by electronic methods.\nOne object of the invention is to create a new model for modelling e.g. the acoustic characteristics of the human speech mechanism, or the producing of speech. Models produced by the method may also be used in speech recognition, in estimating the parameters of a genuine speech signal and in so-called Vocoder apparatus, in which speech messages are transferred with the aid of speech signal analysis and synthesis with a minor amount of information e.g. over a low information rate channel, at the same time endeavouring to maintain the highest possible level of speech quality and intelligibility.\nSince the model of the invention is intended to be suitable for the modelling of events taking place in an acoustic tube in general, the invention is also applicable to electronic music synthesizers.\nThe methods of prior art serving the artificial producing of speech are divisible into two main groups. By the methods of the first group only such speech messages can be produced which have at some earlier time been analyzed, encoded and recorded from corresponding genuine speech productions. Best known among these procedures are PCM (Pulse Code Modulation), DPCM (Differential Pulse Code Modulation), DM (Delta Modulation) and ADPCM (Adaptive Differential Pulse Code Modulation). A feature common to these methods of prior art is that they are closely associated with signal theory and with the general signal processing methods worked out on its basis and therefore imply no detailed knowledge of the character or mode of generation of the speech signal.\nThe second group consists of those methods of prior art in which no genuine speech signal has been recorded, neither as such or in coded form, instead of which the speech is generated by the aid of apparatus modelling the functions of the human speech mechanism. First, from genuine speech are analyzed its recurrent and comparatively invariant elements, phonetic units or phonemes and variants thereof, or phoneme variants, in varying phonetic environments. In the speech synthesizing step, the electronic counterpart of the human speech system, which is referred to as a terminal analog, is so controlled that phonemes and combinations of phonemes equivalent to genuine speech can be formed. To date, these are the only methods by which it has been possible to produce synthetic speech from unrestricted text.\nIn the territory between the said two groups of methods of prior art is located Linear Predictive Coding, LPC, /1/ J. D. Markel, A. H. Gray Jr.: Linear Prediction of Speech, New York, Springer-Verlag 1976. Differing from other coding methods, this procedure necessitates utilization of a model of speech producing. The starting assumption in linear prediction is that the speech signal is produced by a linear system, to its input being supplied a regular succession of pulses for sonant and a random succession of pulses for unvoiced speech sounds. It is usual to employ as transfer function to be identified, an all-pole model (cf. cascade model). With the aid of speech signal analysis, estimates are calculable for the coefficients (a.sub.i) in the denominator polynomial of the transfer function. The higher the degree of this polynomial (which is also the degree of the prediction), the higher is the precision with which the speech signal can be provided with the aid of the coefficient a.sub.i.\nThe filter coefficients a.sub.i are however nonperspicuous from the phonetic point of view. To realize a digital filter using these coefficients is also problematic, for instance in view of the filter hardware structures and of stability considerations. It is partly owing to these reasons that one has begun in linear predicting to use a lattice filter having a corresponding transfer function but provided with a different inner structure and using coefficients of different type.\nIn a lattice filter of prior art, bidirectionally acting and structurally identical elements are connected in cascade. With certain preconditions, this filter type can be made to correspond to the transfer line model of a sound channel composed of homogeneous tubes with equal length. The filter coefficients b.sub.i will then correspond to the coefficients of reflection (.vertline.b.sub.i .vertline.&lt;1). The coefficients b.sub.i are determinable from the speech signal by means of the so-called PARCOR (Partial Correlation) method. Even though the coefficients of reflection b.sub.i are more closely associated with speech production, i.e., with the articulatory aspect, generation of these coefficients by regular synthesis principles has also turned out to be difficult.\nIt is thus understood that speech synthesis apparatus of the terminal analog type, known in prior art, implies that speech production is modelled starting out from an acoustic-phonetic basis. For the acoustic phonation system, consisting of larynx, pharynx and oral and nasal cavities, an electronic counterpart has to be found of which the transfer function conforms to the transfer function of the acoustic system in all and any enunciating situations. Such a time-variant filter is referred to as a terminal analog because its overall transfer function from input to output, or between the terminals, aims at analogy with the corresponding acoustic transfer function of the human phonation system. The central component of the terminal analog is called the sound channel model. As known, this is in use e.g. in vowel sounds and partly also when synthesizing other sounds, depending on the type of model that is being used.\nSince the human phonation system is extremely complex of its acoustical properties, a number of simplifications and approximations must be made when formulating models for practical applications. A problem of principle which figures centrally in such model formulation is that the sound channel is a subdivided system with an acoustic transfer function composed of transcendental functions. Creation of a corresponding terminal analog arrangement using lumped electrical components requires that the acoustic transfer function can be approximated with the aid of rational, meromorphic functions.\nAnother centrally important point is the controllability of the model, that is the number and type of control parameters required in the model to the purpose of creating speech, and the degree in which the group of control parameters meets the requirements of optimal, \"orthogonal\" and phonetically clear-cut selection.\nAs known in the prior art, in constructing sound channel models, the acoustic sound channel is simplified by assuming it to be a straight homogeneous tube, and for this the transfer line equations are calculated (cf. /2/ G. Fant: Acoustic Theory of Speech Production, the Hague, Mouton 1970, Chapters 1.2 and 1.3; and /3/ J. L. Flanagan: Speech Analysis Synthesis and Perception, Berlin, Springer-Verlag 1972, p. 214-228). The assumption is made that the tube has low losses and is closed at one end; the glottis, or the opening between the vocal cords, closed; and the other end opening into the free field. The acoustic load at the mouth opening may be simply modelled either by a short circuit or by a finite impedance Z.sub.r. The acoustic transfer function that is being approximated will then have the form: ##EQU2## where\ny (s)=.alpha.+j.beta.=propagation coefficient\n.alpha.=attenuation factor\n.beta.=.omega./c=phase factor\n.omega.=angular frequency\nc=velocity of sound\nZ.sub.r =radiation load impedance\nZ.sub.o =characteristic impedance of the channel\nl=length of the channel.\nAssuming that the losses of the channel are minor and that the channel terminates in short circuit (Z.sub.r =0), or that the channel is lossless and Z.sub.r is resistive, Equation (1) becomes: ##EQU3## where A, a and k are real. The logarithmic amplitude graph of the absolute value of the transfer function H.sub.A (.omega.) is shown in FIG. 7. The homogeneous sound channel chosen as starting point for the approximation is most nearly equivalent to the situation encountered when pronouncing a neutral vowel (.omega.). The profile of the sound channel and its transfer function are altered for other vowel sounds."}
{"text": "In various communication systems, digital-to-analog converters are used to convert digital signals to analog signals before transmission. Digital-to-analog converters may introduce quantization noise into the analog signals—particularly when a large number of signal levels are used. Examples of techniques that utilize a large number of output levels include Tomlinson-Harashima-Precode and advance modulation schemes such as OFDM and discreet multi-tone modulation.\nTypically, to reduce the effect of quantization noise on system performance, a power spectrum density (PSD) level of the quantization noise should be below a predetermined PSD level of unavoidable noises. A typical requirement is for the quantization noise to have a PSD that is 10 decibels below the PSD of unavoidable noises. Examples of unavoidable noises include additive white Gaussian noise (AWGN), alien cross-talk from other cables or transmitters and quantization noise of the analog to digital converter at the receiver.\nConventional digital-to-analog converter designs produce a quantization noise with a white PSD evenly distributed among all frequency components. However, the communication system performance is often limited by a worst case channel. The frequency response of this channel varies significantly within the transmission bandwidth. As a result, the quantization noise from the transmitter digital-to-analog converter may be shaped by the channel and observed by the receiver.\nThe peak of the quantization noise PSD (shaped by the channel) observed by the receiver must be lower than other noises by a predetermined level. As a result, a large number of bits may be required for the digital-to-analog converter input and the digital-to-analog converter size and complexity are increased. Reducing the size and complexity of a digital-to-analog converter would lower the overall cost of the system.\nReferring now to FIG. 1, a transmitter 8 having an input 10 from an advance modulation scheme or a pre-coding scheme is illustrated. The input 10 generates an N-bit digital input to a truncation module 12. The truncation module 12 truncates the N-bit signal to an M-bit signal, where M is an integer less than N. The truncation module 12 eliminates the least significant bits from the N-bit digital input signal. The M-bit signal is provided to a digital-to-analog converter 14 where it is converted to an analog output signal corresponding to the M-bit signal.\nReferring now to FIG. 2, a signal model illustrating the input signal an, which corresponds to the output of the truncation module 12, is summed with truncation noise qn at a summing module 16. The truncation noise qn is inherent in the truncation module 12. The truncation noise is sometimes referred to as quantization noise.\nReferring now to FIG. 3, a 10GBASE-T transmitter 20 having a pre-coder 18 is illustrated. An input signal ak is provided to a summing module 22, the addition module 22 generates a summed signal dk as will be described below. The signal dk is provided to a modulo operation module 24 where it is converted to a signal sk that has N-bits. Feedback of the signal sk is provided through a feedback filter 26 having a transfer function P(z). The output of the feedback filter 26 is provided to the addition module 22. Referring back to modulo operation module 24, the N-bit signal, sk is provided to the truncation module 28, which truncates the signal to an M-bit signal that is provided to the digital-to-analog converter 32 conversion to an analog signal. The digital-to-analog converter 32 illustrated in FIG. 3 may be implemented with Tomlinson-Harashima-Precode (THP). The approach illustrated in FIGS. 1-3 has quantization noise problems that degrade the performance of the communication system."}
{"text": "Cardiovascular diseases are a class of diseases that involve the heart or blood vessels. Mendis et al., Global Atlas on Cardiovascular Disease Prevention and Control—World Health Organization, World Heart Federation, and World Stroke Organization 2011. Cardiovascular diseases include coronary artery diseases such as angina and myocardial infarction (commonly known as heart attack), stroke, hypertensive heart disease, rheumatic heart disease, cardiomyopathy, cardiac arrhythmia, congenital heart disease, valvular heart disease, carditis, aortic aneurysms, peripheral artery disease, and venous thrombosis. Id.; Lancet 2014, 385, 117-171. Cardiovascular diseases are the leading cause of death globally. Mendis et al., Global Atlas on Cardiovascular Disease Prevention and Control—World Health Organization, World Heart Federation, and World Stroke Organization 2011. In 2013, cardiovascular diseases resulted in 17.3 million deaths (31.5%), up from 12.3 million (25.8%) in 1990. Lancet 2014, 385, 117-171. In the United States, 11% of people between 20 and 40 have a cardiovascular disease, while 37% between 40 and 60, 71% of people between 60 and 80, and 85% of people over 80 have a cardiovascular disease. Go et al., Circulation 2013, 127, e6-e245.\nOne form of cardiac arrhythmia is bradycardia, a slow heart rate condition, which can cause fainting, dizziness, malaise, general weakness, excessive fatigue, chest pain, or failing memory. With no approved drug therapy available for effective treatment of chronic bradycardia, a cardiac pacemaker must often be installed into a patient to sustain his/her life. Although several studies had investigated potential therapy with some drugs (Alboni et al., Am. J. Cardiol. 1990, 65, 1037-1039; Ling et al., Ann. Pharmacother. 1998, 32, 837-839; Benditt et al., Am. J. Cardiol. 1983, 52, 1223-229), the adverse side effects at therapeutic doses prevent routine, long-term use of these previously tested drug candidates. Thus, currently, there is no approved drug worldwide for treating bradycardia. The treatment of choice for chronic, symptomatic bradycardia is limited to the implantation of a cardiac pacemaker. It is estimated that the United States and Europe, as well as other advanced countries, have a rate of pacemaker implantation of about 50-60/100,000 people per annum. A worldwide survey published in 2009 indicates that approximately one million pacemaker implantations (approximately 740,000 new implants) were performed in the 61 countries responding to the survey. By 2028, it is estimated that over 700,000 new pacemaker implantations will be done in the US alone. The 10-year prevalence (US 1999-2008) of clinically defined bradycardia (abnormally slow heart rate, RPR (resting pulse rate)<60 beats/min) is 15.2% for male adults and 6.9% for females. US National Health Statistics Reports 2011, 41, 1-16. Since the human pulse rate is inversely associated with age and the aged population is increasing worldwide, there will be a continually increasing number of patients with bradycardia requiring pacemaker implants. Furthermore, due to the adverse side effects and high cost of pacemakers, many bradycardia patients either cannot or elect not to have a pacemaker implanted even when needed from a medical perspective.\nAs noted, current medical treatment for bradycardia requires a surgical insertion of a cardiac pacemaker. The first use of a buried cardiac pacemaker occurred in Sweden in 1958. Since then, this approach for the treatment of abnormal cardiac rhythms has gradually spread all over the world. In recent years, the quality of pacemakers has improved; and as a result, the use of pacemakers to treat bradycardia has greatly increased. Unfortunately, cardiac pacemakers have some limitations. For example, pacemakers require a surgical procedure for implantation in the human body, and infections requiring device removal do occur with a finite frequency. In addition, pacemakers may not provide a normal physiological heart rate response to exertion or a normal contractile pattern. Because of inadequate or inappropriate heart rate response during exertion, patients may develop significant symptoms related to inadequate cardiac output. In addition, pacing the right ventricle only (as occurs with single or dual chamber pacemakers) may lead to significant left ventricular dysfunction and worsening of heart failure in some patients. Left ventricular failure associated with right ventricular pacing may be ameliorated by implantation of a bi-ventricular pacemaker, but this procedure is time consuming, difficult, and not always possible due to technical issues. As has been recently observed, pacemakers and pacemaker leads may fail and are subject to periodic recalls by sovereign regulatory agencies such as the FDA. Finally, the costs of pacemaker implantation, replacement and follow up are high. Nonetheless, in the absence of a reasonable alternative treatment for bradycardia, doctors and patients usually opt to use a pacemaker when it is necessary. Furthermore, many bradycardia patients delay or forgo the pacemaker surgery due to the potential complications discussed herein and choose to live with bradycardia and the associated potentially fatal medical risks instead. Therefore, there is an unmet need for an effective drug therapy for treating bradycardia."}
{"text": "For a bistable nematic liquid crystal display, state transitions (ON-to-OFF or OFF-to-ON) occur under the influence of appropriately applied dynamic electric fields. See, for example, G. Boyd et al., \"Liquid-Crystal Orientational Bistability and Nematic Storage Effects,\" Appl. Phys. Lett. 36, pp. 556-558 (1980) and J. Cheng et al., \"Threshold and Switching Characteristics of a Bistable Nematic Liquid-Crystal Storage Display,\" Appl. Phys. Lett. 37, pp. 1072-1074 (1981). In particular, vertical electric fields cause orientational director transitions to a predominantly vertical alignment configuration, an ON state, for example, for the liquid crystal molecules. Similarly, horizontal electric fields cause orientational director transitions to a predominantly horizontal alignment configuration, an OFF state, for example, which is topologically distinct from the vertical alignment configuration. See U.S. patent application, Ser. No. 98,976 filed on Nov. 30, 1979, U.S. Pat. No. 4,333,708.\nIn this bistable nematic liquid crystal medium, the horizontal and vertical electric fields are produced with an array of interdigital electrodes as a matrix addressing arrangement for activating and deactivating individual display cells. Although this type of arrangement can provide low to moderate addressing speeds for a moderate electric field strength, it is incapable, from a practical viewpoint, of providing high speed operation. In addition, interdigital electrodes require both a large number of connections per display cell and complex control circuits to activate particular electrodes for switching.\nInterdigital electrodes generate nonhomogeneous electric fields and exhibit two distinct switching thresholds for liquid crystal displays, namely, a longitudinal threshold and a transverse threshold. These thresholds represent minimum electric field strengths necessary for detaching disclinations in the liquid crystal medium from boundary discontinuities in surface alignment or topography. Existence of the two separate thresholds and the creation of nonhomogeneous electric fields substantially impair the effectiveness of interdigital electrodes for high speed matrix addressing purposes in liquid crystal displays.\nHomogeneous or uniform vertical electric fields can be generated by an array of orthogonally disposed, continuous uniform strip electrodes as the matrix addressing arrangement. This type of electrode arrangement facilitates transitions from the horizontal state to the vertical state, that is, OFF to ON state transitions. Furthermore, it exhibits a single sharp switching threshold for horizontal to vertical state transitions which is of sufficient magnitude and definition to assure reliable, high speed switching. However, the array of strip electrodes is unsuited for generating the fields necessary for vertical to horizontal state transitions, that is, ON to OFF state transitions."}
{"text": "Lighting systems may be deployed to provide lighting for various environments such as parking lots, roadways, sidewalks, structures, etc. In some of these environments, the illumination from the fixtures may be managed by photocontrols that adjust the amount of light produced based upon the amount of ambient light detected by the photocontrol. For example, a photocontrol placed outside may switch off light output from the light fixture during daylight hours and switch on the light output during the evening.\nHowever, changing the programming of these photocontrols to, for example, disable light output or temporarily switch on the light, may be a complicated task. To accomplish the task, a user often needs to have physical access to the photocontrol, which may be installed in a location that is difficult to access, such as a atop a 40-ft. light pole. The difficulty can be magnified when there are many photocontrols which need to be adjusted."}
{"text": "The present invention relates to intermediate molecular weight shaped polyethylene articles such as polyethylene fibers exhibiting relatively high tenacity, modulus and toughness, and to products made therefrom. The polyethylene article is made by a process which includes the step of stretching a solution of polyethylene dissolved in a solvent at a stretch ratio of at least about 3:1.\nPolyethylene fibers, films and tapes are old in the art. An early patent on this subject appeared in 1937 (G.B. No. 472,051). However, until recently, the tensile properties of such products have been generally unremarkable as compared to competitive materials, such as the polyamides and polyethylene terephthalate. Recently, several methods have been discovered for preparing continuous low and intermediate molecular weight polyethylene fibers of moderate tensile properties. Processes for the production of relatively low molecular weight fibers (a maximum weight average molecular weight, Mw, of about 200,000 or less) have been described in U.S. Pat. Nos. 4,276,348 and 4,228,118 to Wu and Black, U.S. Pat. Nos. 3,962,205, 4,254,072, 4,287,149 and 4,415,522 to Ward and Cappaccio, and U.S. Pat. No. 3,048,465 to Jurgeleit. U.S. Pat. No. 4,268,470 to Cappaccio and Ward describes a process for producing intermediate molecular weight polyolefin fibers (minimum molecular weight of about 300,000).\nThe preparation of high strength, high modulus polyolefin fibers by solution spinning has been described in numerous recent publications and patents. German Off. No. 3,004,699 to Smith et al. (Aug. 21, 1980) describes a process in which polyethylene is first dissolved in a volatile solvent, the solution is spun and cooled to form a gel filament, and, finally, the gel filament is simultaneously stretched and dried to form the desired fiber. U.K. Patent Application No. 2,051,667 to P. Smith and P. J. Lemstra (Jan. 21, 1981) discloses a process in which a solution of a polymer is spun and the filaments are drawn at a stretch ratio which is related to the polymer molecular weight, at a drawing temperature such that at the draw ratio used, the modulus of the filaments is at least 20 GPa (the application notes that to obtain the high modulus values required, drawing must be performed below the melting point of the polyethylene; in general, at most 135.degree. C.). Kalb and Pennings in Polymer Bulletin, Volume 1, pp. 879-80 (1979), J. Mat. Sci., Vol. 15, pp. 2584-90 (1980) and Smook et al. in Polymer Mol., Vol 2, pp. 775-83 (1980) describe a process in which the polyethylene is dissolved in a non-volatile solvent (paraffin oil), the solution is cooled to room temperature to form a gel which is cut into pieces, fed to an extruder and spun into a gel filament, the gel filament being extracted with hexane to remove the parafin oil, vacuum dried and stretched to form the desired fiber.\nMost recently, ultra high molecular weight fibers have been disclosed. U.S. Pat. No. 4,413,110 to Kavesh and Prevorsek describes a solution spun fiber of from 500,000 molecular weight to about 8,000,000 molecular weight which exhibits exceptional modulus and tenacity. U.S. Pat. Nos. 4,430,383 and 4,422,993 to Smith and Lemstra also describe a solution spun and drawn fibers having a minimum molecular weight of about 800,000. U.S. Pat. No. 4,436,689 to Smith, Lemstra, Kirschbaum and Pijers describes solution spun filaments of molecular weight greater than 400,000 (and an Mw/Mn&lt;5). In addition, U.S. Pat. No. 4,268,470 to Ward and Cappacio also discloses high molecular weight polyolefin fibers.\nIn general, the known processes for forming polyethylene and other polyolefin fibers may be observed as belonging in one of two groups: those which describe fibers of low average molecular weight (200,000 or less) and those which describe fibers of high average molecular weight (800,000 or more). Between the two groups, there is a molecular weight range which has not been accessible to the prior art methods for preparing fibers of high tensile properties.\nThere are advantages to the molecular weight ranges thus far mastered. Lower molecular weight polymers are generally synthesized and processed into fibers more easily and economically than high molecular weight fibers. On the other hand, fibers spun from high molecular weight polymers may possess high tensile properties, low creep, and high melting point. A need exists for fibers and methods which bridge this gap, combining good economy with moderate to high tensile properties. Surprisingly, our process makes it possible to accomplish these results."}
{"text": "For example, a weatherstrip is mounted on an opening edge of a vehicular body side opening section opened or closed by means of a vehicular trunk in order to secure a sealing characteristic. This weatherstrip 30 includes a welt section 32 in a substantially letter U shape in cross section and fitted to a flange 31 mounted on a peripheral edge section of the vehicular body side opening section and a trim lip 34 extended from an outer surface of this welt section 32 and which covers and hides an end edge of vehicular body trim 33. This trim lip 34 causes the end edge of vehicular body trim 33 not to be exposed and an outer appearance quality is improved. Welt section 32 is provided with holding lips 35, 35 extended toward obliquely a summit section from an inner surface of the respective side wall sections mutually connected from the summit section. When these holding lips 35 are elastically contacted on flange 31, welt section 32 is held by flange 31 (for example, refer to a Patent document 1)."}
{"text": "1. Field of Our Invention\nOur invention relates to apparatus for testing mechanical heart valves, and in particular mechanical heart valves having a rigid annulus and one or more pivoting leaflets mounted therein.\n2. Description of Related Art\nVarious types of heart valve prostheses have been proposed, and many give generally satisfactory operation. One popular design for a heart valve prosthesis includes an annular valve body in which a pair of opposed leaflet occluders are pivotly mounted. The occluders are moveable between a closed, mated position, blocking blood flow in an upstream direction, thereby minimizing regurgitation, and an open position, allowing blood flow in a downstream direction. One such heart valve is described, for example, in U.S. Pat. No. 5,147,390 to Campbell, which patent is assigned to CarboMedics, Inc., the assignee of our invention.\nA mechanical heart valve, such as that described in the '390 patent, can be expected to open and close a great number of times during its use. It is desirable to minimize, in so far as possible, the number of failures experienced in the use of a prosthetic heart valve. Testing for function is therefore an important part of prosthetic heart valve development and manufacture. Heart valve function testers are known which open and close the mechanical heart valve in an in vitro environment, mimicking the action of the heart. Fluid is forced past the valve to open the valve. An existing back pressure is then allowed to close the valve when the pulsatile forward pressure is removed.\nThere is, however, another possibility for defects which can be tested. Minute cracks or other surface defects in pivots of leaflet occluders or in pivot recesses of the valve body are difficult to detect. It is known, however, from the application of fracture mechanics, that cracks below a certain maximum size will not cause failure. Cracks or other surface defects larger than the maximum allowable size can be detected by applying a proof test load to the component. The proof test load should be some multiple of the functional vivo load, to provide a factor of safety associated with the test. The primary purpose of proof testing is to detect components of heart valves with flaws larger than a specific size."}
{"text": "Nitrogen is an important plant nutrient. In addition to phosphorous, potassium, and other nutrients, nitrogen is needed to support the growth and development of plant life. Some plants, such as legumes, through a symbiotic relationship with Rhizobium bacteria, fix elemental nitrogen from the atmosphere and fix this nitrogen into the soil. However, most plants grown to produce human and animal food require the use of nitrogen fertilizer in order to sustain their agricultural production.\nThe most widely used and agriculturally important high-analysis nitrogen fertilizer is urea, CO(NH2)2. While most of the urea currently produced is used as a fertilizer in its granular form, urea-based fluid fertilizers are also well known. As used herein, the term “fluid fertilizers” encompasses liquid fertilizers, i.e. aqueous solutions of fertilizers, and suspension fertilizers, i.e. fertilizer compositions which in addition to water and water-soluble components also contain insoluble components kept in suspension by a suspending agent, such as clay. Suspension fertilizers are excellent carriers for pesticides and micronutrients.\nThe most commonly known urea-based liquid fertilizer is an aqueous solution of urea and ammonium nitrate, referred to in the fertilizer trade as UAN solution. These fluid fertilizers are used on a variety of crops, such as corn and wheat. When applied to moist soil, the urea component of the fluid fertilizer becomes a source of ammonia as a result of hydrolysis catalyzed by urease, an enzyme produced by numerous fungi and bacteria. The process is fully disclosed in U.S. Pat. No. 5,364,438, which is hereby incorporated by reference. Unfortunately, the urease-catalyzed hydrolysis often converts the urea to ammonia more quickly than it can be absorbed by the soil, resulting in undesirable ammonia loss to the atmosphere through a process called volatilization.\nUrease inhibitors slow down the conversion of urea to ammonia, extending the period of nitrogen release. One urease inhibitor is N-(n-butyl) thiophosphoric triamide (NBPT). When incorporated into a urea-based fertilizer, NBPT reduces the rate at which urea is hydrolyzed to ammonia. This allows the nutrient nitrogen to be available to the soil and plants over a longer period of time. NBPT is widely recognized as an effective urease inhibitor, but NBPT is notoriously difficult to handle, because industrial grade NBPT is a waxy, sticky, heat-sensitive and water-sensitive material. As a result, a solvent system was disclosed in U.S. Pat. No. 5,364,438, that allowed the solution of NBPT to be applied to UAN. However, this solution had stability problems as well as problems in delivering and metering into fluid fertilizer.\nU.S. Pat. No. 5,352,265, which is hereby incorporated by reference, discloses a granular fertilizer comprising about 90 to 99% urea, 0.02 to 0.5% NBPT and about 0 to 2.2% DCD. The NBPT is added to molten urea as a concentrated solution in an amide solvent. The DCD is added to the urea melt as a solid. This granular product is made to apply directly to the field crop.\nThe present invention of a dry flowable additive, is prepared by coating a solid urea-formaldehyde polymer (UFP) with a solution or suspension of NBPT in a liquid solvent, preferably an amide solvent. Optionally, the coated UFP may be blended with solid DCD. Prior to application of the fertilizer to the field crop, the dry flowable additive is blended with a UAN solution or aqueous urea, to form the fluid urea-containing fertilizer composition, or blended with solid or molten urea to form a solid urea-based fertilizer. The present invention provides a fluid or solid fertilizer composition that is easy to handle and stable when stored."}
{"text": "Pyrolytic incineration operates on a starved air principle, and air supplied to the pyrolysis chamber generally constitutes less than half of the stoichiometric air requirement for combustion of the waste. The low air supply rate achieves partial combustion and vaporization of the waste, and results in low gas velocity and turbulence in the pyrolysis chamber which minimizes mechanical entrapment of particulate matter in the waste gases.\nThe waste gases from the pyrolysis chamber, then pass into the thermal reactor section of the incinerator which is located in the stack. Atmospheric air is drawn into the thermal reactor section to achieve substantially complete combustion of the combustible waste gases in a second zone of combustion.\nThe normal incinerator is designed for peak thermal capacities, and the peak capacity is rarely attained except when waste material is being charged into the combustion chamber. Consequently, the incinerator is normally operating well below peak capacity.\nCertain liquid wastes, such as solvents, paint, lacquer, and the like, have a high BTU content, while other water base liquid wastes have a relatively low BTU content. In addition, the viscosity and solids content of the liquid waste can vary considerably.\nIn the past, incinerators have been designed to separately burn either solid or liquid waste materials and in some cases, incinerators have been designed with sequential combustion chambers to handle both liquid and solid waste. However, none of the incinerators, as used in the past to handle both solid and liquid waste materials have been programmed to utilize the spare, below-peak capacity of the solid waste incinerator."}
{"text": "Currently, computing devices, particularly those associated with machine learning for autonomous and semi-autonomous vehicles, store three dimensional (3D) models that are used as training data for machine learning. This type of training data is computationally expensive and requires a large amount of storage space. As such, it may be difficult for on-board vehicle processors to effectively utilize the training data."}
{"text": "1. Field of the Invention\nThe present invention relates to high density memory devices based on phase change memory materials, including chalcogenide based materials and on other programmable resistance materials, and methods for manufacturing such devices.\n2. Description of Related Art\nPhase change based memory materials, like chalcogenide based materials and similar materials, can be caused to change phase between an amorphous state and a crystalline state by application of electrical current at levels suitable for implementation in integrated circuits. The generally amorphous state is characterized by higher electrical resistivity than the generally crystalline state, which can be readily sensed to indicate data. These properties have generated interest in using programmable resistance material to form nonvolatile memory circuits, which can be read and written with random access.\nThe change from the amorphous to the crystalline state is generally a lower current operation. The change from crystalline to amorphous, referred to as reset herein, is generally a higher current operation, which includes a short high current density pulse to melt or breakdown the crystalline structure, after which the phase change material cools quickly, quenching the molten phase change material and allowing at least a portion of the phase change material to stabilize in the amorphous state.\nThe magnitude of the current needed for reset can be reduced by reducing the size of the phase change material element in the cell and/or the contact area between electrodes and the phase change material, so that higher current densities are achieved with small absolute current values through the phase change material.\nOne approach to reducing the size of the phase change element in a memory cell is to form small phase change elements by etching a layer of phase change material. However, reducing the size of the phase change element by etching can result in damage to the phase change material due to non-uniform reactivity with the etchants which can cause the formation of voids, compositional and bonding variations, and the formation of nonvolatile by-products. This damage can result in variations in shape and uniformity of the phase change elements across an array of memory cells, resulting in electrical and mechanical performance issues for the cell.\nAdditionally, it is desirable to reduce the cross-sectional area or footprint of individual memory cells in an array of memory cells in order to achieve higher density memory devices. However, traditional field effect transistor access devices are horizontal structures having a horizontally oriented gate overlying a horizontally oriented channel region, resulting in the field effect transistors having a relatively large cross-sectional area which limits the density of the array. Attempts at reducing the cross-sectional area of horizontally oriented field effect transistors can result in issues in obtaining the current needed to induce phase change because of the relatively low current drive of field effect transistors.\nThus, memory devices including both vertically and horizontally oriented field effect transistors have been proposed. See, for example, U.S. Pat. No. 7,116,593. However, the integration of both vertically and horizontally oriented field effect transistors can be difficult and increase the complexity of designs and manufacturing processes. Thus, issues that devices having both vertically and horizontally oriented field effect transistors need to address include cost and simplicity of manufacturing.\nAlthough bipolar junction transistors and diodes can provide a larger current drive than field effect transistors, it can be difficult to control the current in the memory cell using a bipolar junction transistor or a diode adequately enough to allow for multi-bit operation. Additionally, the integration of bipolar junction transistors with CMOS periphery circuitry is difficult and may result in highly complex designs and manufacturing processes.\nIt is therefore desirable to provide both vertically and horizontally oriented field effect transistors on the same substrate that are readily manufactured for use in high-density memory devices, as well as in other devices that may have a need for both types of transistors on one chip. It is also desirable to provide memory devices providing the current necessary to induce phase change, as well as addressing the etching damage problems described above."}
{"text": "1. Field of the Invention\nThe present invention relates to virtual reality displays.\n2. Related Art\nPrior methods of virtual reality display systems deployed head or helmet mounted display that placed a viewing screen directly in front of the user's eyes and recorded the movement of the users head to determine what should be shown on the display. Thus, when the head turned to one side, the display was refreshed to show what was in the virtual world in the direction they turned their head."}
{"text": "1. Technical Field\nThe present invention relates to an improved information-retrieval apparatus. In particular, the present invention relates to an improved digitally-based information-retrieval apparatus. More particularly, the present invention relates to magnetic storage media, such as digital recording tapes. Still more particularly, the present invention relates to magnetic recording heads for writing and reading data to digital recording tapes.\n2. Description of the Related Art\nVarious magnetic recording techniques exist for recording data to and from magnetic storage media, such as magnetic tape. Magnetic tapes are used for data storage in computer systems requiring data removability, low-cost data storage, high data-rate capability, high volumetric efficiency and reusability. The constantly increasing operational speeds of digital computers are creating a demand for corresponding increases in the data storage capacities of magnetic tape recording and reproducing systems, while maintaining the special requirements of high speed digital tape systems.\nTape recording and reproducing systems for use as computer data storage devices are required to provide high data transfer rates and to perform a read check on all written data. To satisfy these requirements, conventional tape systems typically employ methods of recording known as linear recording, in which the tracks of data lie parallel to each other and to the edge of the tape. Linear recording techniques offer high data transfer rates. However, it is desirable to obtain even higher data densities while retaining the advantages of such recording techniques.\nDigital linear tape (DLT) is a magnetic linear tape medium that is increasingly being utilized as a medium for data storage. DLT is a magnetic storage medium used to back up data, typically in computer systems. DLT allows for the rapid transfer of data, in comparison to other tape storage technologies. For example, various forms of magnetic read/write heads can be utilized in association with servo mechanisms to read and write data to and from a track of a particular DLT.\nBecause DLT is currently being utilized as an important tool for data storage, it is desirable to increase the recording density, thus allowing for the faster and more efficient retrieval and writing of data. One method of increasing this storage density involves azimuth recording. The term \"azimuth\" refers to the horizontal angular distance from a particular reference direction. The use of the word \"azimuth\" in \"azimuth recording\" thus suggests a form of angular recording.\nAzimuth recording involves the use of a rotating recording head, such that data tracks on a tape may be recorded at different angles with respect to the edge of the tape. Azimuth recording results in a recorded track pattern in which the magnetization directions of adjacent data tracks lie at different azimuth angle to each other. To date, most recording systems have relied strictly on magnetic heads which contain read/write elements but which record only vertically, thus not allowing for angular or \"azimuthal\" recording of data. One of the principal advantages of azimuth recording over non-azimuth recording is that azimuth recording promotes very high data track packing. Azimuth recording provides much denser track packing than regular track packing spacing because regular track packing spacing typically requires gaps between tracks.\nThose systems which do attempt to implement azimuth recording techniques are faced with the challenge of providing fine positioning servo tracking. Servo tracking techniques have been developed to reduce the effects of tracking error and thus improve the data capacity of tape systems. Known servo techniques vary widely, but most involve methods of dynamically moving a read head gap to continually reposition it over a written servo track. The movement of a servo read head gap compensates for lateral tape motion during a read. However, lateral tape motion during writing is not controlled with respect to the write head gap. Thus, the distance between tracks is still limited to the magnitude of the lateral tape motion in order to avoid over-writing previously written tracks.\nBased on the foregoing it can be appreciated that a need exists for an improved azimuth recording system which does not encounter problems associated fine positioning servo tracking. A need also exists for an inexpensive and easy to implement apparatus and method which provides fine positioning servo tracking. It is believed that the apparatus and method presented herein solves these problems."}
{"text": "(a) Field of the Invention\nThe present invention relates to a liquid crystal display and a drive method thereof, in which liquid crystal molecules respond fast even at data voltages of an intermediate grayscale level. More particularly, the present invention relates to a liquid crystal display and a drive method thereof, which improves a liquid crystal response speed with respect to the application of a gate voltage of a twisted nematic liquid crystal display.\n(b) Description of the Related Art\nA twisted nematic liquid crystal display (TN LCD) has the advantages of enabling control at very thin profile configurations, and of consuming very little power. However, the drawbacks of TN LCDs are that they have slow response speeds with respect to applied voltages, and a limited viewing angle.\nFIG. 1 shows a graph of response curves when a voltage is applied to pixels of a TN LCD.\nAs shown in FIG. 1, a response time of twisted nematic liquid crystals is roughly 15–17 ms from the moment a voltage is applied, and when the applied voltage is switched off, a response time of approximately 20 ms is required. Accordingly, it is difficult to realize images containing a large amount of data.\nVarious configurations are used to improve response speeds. These include the surface stabilized ferroelectric liquid crystal display (SSFLCD) and the anti-ferroelectric liquid crystal display (AFLCD). However, in these LCDs, alignment and the display of grayscale levels are difficult to obtain, and a high reset voltage is required such that practical applications of the LCDs are not fully feasible.\nIn more traditional configurations, the slow response speeds make the display of certain images (e.g., moving images) unclear since these images require the display of large amounts of grayscale levels during a short interval of time. Therefore, the TN LCD particularly needs an improvement in response speeds."}
{"text": "Reduced engine operation times in plug-in hybrid electric vehicles (PHEVs) enable fuel economy and reduced fuel emissions benefits. However, the shorter engine operation times can lead to longer refueling intervals ultimately resulting in fuel in a fuel tank of the vehicle becoming old, or sour. Fuel souring may cause acid formation and/or waxing, for example.\nOne approach to address potential souring of on-board fuel is to force engine on operation to a greater extent, even if not needed. However, the inventors herein have recognized a problem with such an approach. Namely, user satisfaction with the plug-in vehicle may become degraded because the user may be aiming to minimize addition of fuel, and forcing engine operation to utilize stale fuel is directly contrary to the user's goal.\nThus, the inventors herein have devised an approach to at least partially address the issue described above. In one example, a method for a vehicle including an engine and a motor includes displaying a recommended engine fuel fill amount based on a history of actual fuel usage. As an example, the recommended engine fuel fill amount may be determined based on an amount of fuel consumed over a duration, such as since a last fuel refill.\nIn this way, forced engine operation may be reduced such that user satisfaction may be improved. For example, by determining the amount of fuel consumed since the last fuel refill, an operator of the vehicle may be informed of his/her fuel usage via a display in the vehicle. As such, when the operator refills a fuel tank of the vehicle, the operator may refill the fuel tank with just enough fuel for a selected duration (e.g., three months) such that fuel souring may be prevented. Further, as the operator follows the recommended fuel fill amount, a frequency of forced engine operation (e.g., a fuel maintenance mode) may be reduced, as there is no longer excess fuel that needs to be consumed before souring can occur.\nIt should be understood that the summary above is provided to introduce in simplified form a selection of concepts that are further described in the detailed description. It is not meant to identify key or essential features of the claimed subject matter, the scope of which is defined uniquely by the claims that follow the detailed description. Furthermore, the claimed subject matter is not limited to implementations that solve any disadvantages noted above or in any part of this disclosure."}
{"text": "Currently available laser delivery techniques used for dermatology, such as for tattoo removal and medical laser ablation, typically operate in the UV, visible, or infrared spectra. Furthermore, current dermatological applications typically involve the propagation of waves using picosecond, nanosecond, microsecond, or millisecond pulse widths, and also involve continuous wave lasers. Often, a fiber optic bundle, articulated mirror system, or the like directs the laser light from a beam source into a hand-piece control unit positioned above the skin and aimed toward the target. This creates an area of free space (i.e., an air gap) between the hand-piece and the target. Because these procedures propagate a laser through free space to illuminate target tissue, they allow for a dangerous degree of electromagnetic energy to be released, both when the wave propagates in free space and when the propagated light reflects off of the target tissue. Under typical circumstances, even a diffuse reflectance of 1% of the transmitted wave into the eyes of an operator or patient is enough to cause permanent ocular damage.\nIn addition, conventional laser systems exhibit poor heat dissipation at the tissue-air interface. This can exacerbate the negative thermal effects at the tissue surface, which result in potentially drastic changes to the tissue being treated, diminishing overall effectiveness of the treatment, prolonging the time needed to recover between treatments, increasing the number of treatments required, etc. For example, excessive heat can damage or permanently scar the topmost layers of skin, reduce the efficacy of subsequent treatments, and even vaporize water within the tissue. In addition to damaging the target tissue, vaporizing water within the tissue increases the relative fat density at the surface, which causes even more back-scattered light to reflect from the tissue surface during treatment."}
{"text": "1. The Field of the Invention\nThis invention relates to computerized methods for testing and tracking and, more particularly, to novel systems and methods for testing, tracking, and correcting errors due to software or hardware.\n2. Background\nProduct development cycles have become shorter and shorter. More of the responsibility for testing and “debugging” products falls to the actual beta testers or alpha testers. Nevertheless, products are continuing their development cycle well into their marketing bases.\nFor example, software is often released for public purchase before the known errors from beta testing have been cured. Hardware is often likewise premature, and more difficult to correct. Alternatively, beta testing may be inadequate, leaving various problems extant within either hardware, software, or a combination thereof.\nPurchasers are often left with a need for identification and cure of errors in commercially available software and hardware. In some instances, product manufacturers and suppliers actively solicit comments, improvements, detection and identification of errors, and the like. In other situations, manufacturers and marketers of products are not so forthcoming.\nFor example, occasionally, problems are comparatively esoteric, and may occur only in a few rare conditions or instances. Nevertheless, some errors occur with sufficient regularity as to seriously encumber users unaware of the existence of such product flaws.\nIn recent years, computer and software manufacturers have been repeatedly surprised, even amazed, at the groundswell of opposition to products that are not adequately tested, supported, corrected, recalled, or otherwise identified as having correctable flaws.\nSoftware, in particular, has arrived at a new threshold of pain for purchasers and users. Never since the advent of government agencies for consumer protection against fraud, product failure, product inadequacy, manufacturer non-responsiveness, and the like, have so many dollars of product value been subject to such massive amounts of owner and operator time in order to obtain the purported benefits of the products.\nSome manufacturers are swift to seek out and post notification of errors existing in their products. Typically, errors are identified, with associated patches for correcting the errors. In some cases, products are recalled. With the advent of the world wide web, a host of users may provide a corresponding host of error corrections, all freely available to users interested in improving the performance or reliability of a purchased software or hardware product in the computer industry.\nHistorically, a manufacturer or other purveyor of a computer-related product may face a dilemma with respect to certain product flaws. To the extent that an error, built into or programmed into a computer-related product, is comparatively esoteric and unlikely to cause problems for the majority of users, a manufacturer or developer may prefer to ignore it. To the extent that such a flaw or error is ubiquitous and likely to cause pervasive and obvious problems, a manufacturer may prefer to cure the problem. Similarly, to the extent that a problem is likely to cause a comparatively small disruption of promised service, a manufacturer may choose to ignore it. Alternatively, to the extent that a problem is likely to cause serious economic damages to a commercial or industrial user of a software product or physical damage to persons or property as a direct result of the failure of a computer-based product, a manufacturer will take appropriate steps to find a correction to the problem, announce the presence of the flaw and the availability of a corrective measure, and seek to bring all copies of the product into compliance with a corrected version thereof.\nNevertheless, product improvement is largely a matter of motivation. Motivation may arise from personal interest, individual or enterprise-wide frustration, desirability of a result, previous experiences and expectations, and the like. In current process for product improvement, little incentive exists to provide for skilled third parties to improve marketed products. By the same token, manufacturers, whether large or small, may have limited motivation, resources, or the like to locate and correct errors. In fact, a certain motivation may exist to not seek out errors, nor to highlight them, nor even to repair them, in many instances.\nWhat is needed is a mechanism, whereby software and hardware products related to computer systems may be improved profitably by third parties. Likewise, what is needed is an apparatus and method for consistently providing the necessary resources for testing, correction, notification, and product redistribution for products and upgrades related to computer-related based products, whether software or hardware."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus, such as an electronic copying apparatus and a laser printer, and more particularly to an image forming apparatus which has a fixing unit for fixing a developer (e.g., powdered toner) transferred onto a sheet.\n2. Description of the Related Art\nIn recent years, an image forming apparatus (e.g., an electronic copying apparatus and a laser printer) has come to employ the so-called clam shell structure, for easy maintenance and for easy removal of sheets causing an abnormal transfer state. In the clam shell structure, the housing is divided into first and second units (i.e., upper and lower units), with the boundary therebetween defined by a sheet transport path formed inside the apparatus, and wherein the first unit can be opened in a direction away from the second unit.\nThe image forming apparatus comprises a fixing means for fixing a developer image to a sheet. Normally, the fixing means is of a heating roller type. It includes a pair of fixing rollers (namely, a heating roller and a pressing roller) which are in rolling contact with each other, and a releasing claw which is used for releasing a sheet from the fixing rollers, to prevent the sheet from being wound around the fixing rollers.\nPublished Unexamined Japanese Patent Application (PUJPA) No. 63-6588 discloses an example of a conventional image forming apparatus which is the clam shell type and which employs a fixing means of the heating roller type. As may be understood from the example in the reference, the releasing claw remains engaged with the fixing rollers even after the first unit is opened or separated from the second unit.\nIn the conventional image forming apparatus, the releasing claw remains engaged with the fixing rollers even after the first unit is opened, as mentioned above. If the releasing claw is touched by something and is thus subject to an external force, an impact is applied between the releasing claw and the fixing rollers. In such a case, it is likely that the tip end of the releasing claw will be broken or the surfaces of the fixing rollers will be scratched, so that a satisfactory fixing operation or a reliable releasing operation cannot be performed thereafter."}
{"text": "1. Field of the Invention\nThis invention relates to a process for the preparation of crystalline oxytitanium phthalocyanine. Particularly, it relates to a novel process for the preparation of crystalline oxytitanium phthalocyanine suitable for use in an electrophotographic photoreceptor.\n2. Description of the Prior Art\nIn this field, it is well known that the oxytitanium phthalocyanine may be produced either in a semi-stable .alpha.-crystal phase or in a stable .beta.-crystal phase depending on the production conditions. Further, it is also known that those two crystal phases may be converted from one to another under the influences of physical distortion, organic solvent and/or heat. For example, .alpha.-phase oxytitanium phthalocyanine can be converted into .beta.-phase by heating in an organic solvent such as N-methylpyrrolidone.\nOn the other hand, Japanese Patent Application Laying Open (KOKAI) No. 63-20365 discloses that crystalline oxytitanium phthalocyanine of which crystal phase can not be classified neither as the .alpha.-phase nor as the .beta.-phase and which shows an intense peak in the X-ray diffraction spectrum at a Bragg angle (2.theta.) of 27.3.degree. can be produced by forming an aqueous suspension of .alpha.-phase oxytitanium phthalocyanine obtained by the acid paste method, adding an aromatic hydrocarbon solvent into the suspension and heating the suspension, and that the obtained phthalocyanine is an useful material particularly as a recording material of an optical disc.\nThe object of the present invention is to provide a novel and industrially applicable process for the production of oxytitanium phthalocyanine having a crystal phase useful for various purposes such as for use in an electrophotographic photoreceptor."}
{"text": "In general, a solar battery device includes a configuration shown in FIG. 1. In FIG. 1, the reference numeral 1 represents a p-type semiconductor substrate formed in a plate-like shape, and having a size ranging from 100 to 150 mm square and a thickness ranging from 0.1 to 0.3 mm. The p-type semiconductor substrate herein includes a polycrystalline or monocrystalline silicon and the like and is doped with a p-type impurity such as boron and the like. A manufacturing method for the solar battery device will be hereinafter described. First, this substrate is doped with an n-type impurity such as phosphorus and the like to form an n-type diffusion layer 2. Next, an antireflection film 3 such as silicon nitride (SiN) and the like is provided. Then, a conductive aluminum paste is printed on a backside of the substrate by a screen printing method. Thereafter, by drying and firing the conductive aluminum paste, a backside electrode 6 and a BSF (Back Surface Field) layer 4 are formed simultaneously. Successively, a conductive silver paste is printed on a front-side of the substrate. Then, the conductive silver paste is dried and fired to form front-side electrodes 5. With regard to the solar battery device manufactured in such a manner, the front-side electrodes 5 include busbar electrodes and current-collecting finger electrodes. The busbar electrodes are for taking out light-generating current generated by the solar battery device to the outside thereof. The current-collecting finger electrodes are connected to these busbar electrodes. Hereinafter, while a surface of the substrate which is to be a light receiving surface side of the solar battery is referred to as a front-side, a surface of the substrate which is to be the opposite side of the light receiving surface is referred to as a backside.\nWith regard to the solar battery device manufactured in such a manner, the electrodes are generally formed by the screen printing method and firing as mentioned above. In the screen printing method, in order to form the finger electrodes and the busbar electrodes on the light receiving surface of the solar battery cell, for example, a conductive paste containing a silver powder, an organic vehicle, and a glass frit is generally used. Solids of various types of inorganic oxides or conductive materials and the like may be added to this conductive paste to improve performance thereof. When applying this conductive paste to a predetermined position of the semiconductor substrate by the screen printing method and firing the paste, the silver powder sinters each other under high-temperature to form silver electrodes. At the same time, the glass frit is softened to melt the antireflection film and reach the n-type diffusion layer, and the silver electrodes are electrically connected to the n-type diffusion layer. In general, such a method is called Fire Through, which is adopted as a method for forming electrodes of various solar battery cells.\nWith regard to the aforementioned method for forming electrodes, in order to fire the electrodes, the semiconductor substrate should be subjected to high-temperature treatment at 600° C. or more. Due to this high-temperature treatment, the semiconductor substrate is damaged by heat. Alternately, a contaminant which is a lifetime killer gettering upon the diffusion layer may be released inside the semiconductor substrate, which decreases a lifetime of the semiconductor substrate. Furthermore, the electrodes formed by Fire Through are obtained by sintering the conductive particle for a short time. Therefore, the following electrodes may be easily formed, which is a problem. That is, electrodes which have small density compared to electrodes formed by plating and have plenty cavities in front-sides or insides of the electrodes, and each area of those electrodes connected to the semiconductor substrate is not uniform and is easily peeled off, for example. Such a decrease in the lifetime and abnormality in the electrodes may cause problems in performance or long-term reliability of the solar battery cell so that solutions to these problems are demanded.\nIn order to solve the problems, for example, Patent Document 1 discloses a method where a solar battery cell formed by firing electrodes is subjected to heat treatment under an atmosphere including at least hydrogen gas to improve contact resistance of the electrodes. However, in the method disclosed in Patent Document 1, a process is added after firing, which leads to increase costs. Furthermore, there is a problem in safety of the process since hydrogen gas, which is difficult to handle, is used. Therefore, a much simpler method to solve the problems is demanded."}
{"text": "This invention relates generally to the field of collapsible reusable shipping containers of the type disclosed in my prior U.S. Pat. No. 3,443,737 granted Apr. 13, 1969; and more particularly to an improved sliding door construction which may be incorporated into a vertical wall of such container to impart greater structural rigidity thereto under fluid loads which tend to exert a hydraulic outward pressure against the inner surfaces of the container.\nIn the container disclosed in the above patent, the door is supported in a through opening in a side wall by a pair of extrusions of generally H-shaped cross section, each of which define opposed recesses engaging a vertically disposed edge of the door opening, and a corresponding vertical edge on the door. While this arrangement does guide the door smoothly during opening and closing movement, the pressure on the edges of the abutting wall is entirely frictional in nature, and when the container is loaded, for example, with a relatively heavy particulate load, resistance to outward deformation of the wall and the door depends almost entirely upon the engagement of the wall at the lower edges thereof with the rigid pallet, and the engagement of the upper edges thereof beneath the rim of the detachable lid. Depending upon the height and thickness of the walls of the container, and the density of the fluid load, this construction has at times proven inadequate. It is, of course, possible to rigidly interconnect the extrusions to the edges of the wall, but since the door edges must remain free to move in order to selectively open the door, this expedient has but limited utility."}
{"text": "1. Field of the Invention\nThe present invention relates to a clock processing part and more particularly to an internal clock generating circuit and a method for the same.\n2. Description of the Related Art\nTypically, a central processing unit (CPU) and a semiconductor memory device are interconnected through a signal bus. In such a case, the CPU and the semiconductor memory device function as master and slave, respectively. The CPU master transmits data including address, command, writing data, and a clock required for sampling data to the memory device slave.\nAn external clock transmitted through the signal bus may be a clock aligned or centered to the data as shown in FIGS. 1 and 2. The slave memory receives the external clock and generates an internal clock needed for sampling data. In order to correctly sample data, the internal clock should be a data-centered clock as shown in FIG. 1b. If an external clock is a data-aligned clock, it is relatively difficult to generate a data centered internal clock. However, a gradual increase in the data rate/pin decreases the number of valid data windows. If the data and clock have slightly different paths in the system, there may be a bigger skew between the clock and pin that will be applied to the slave. The problem gets worse in a double data rate (DDR) product, which receives two pieces of data at one clock cycle as shown in FIG. 2b, compared to a single data rate (SDR) product as shown in FIG. 2a. \nWhen an external clock is centered or aligned to data, a system designer desires to adopt to the slave a function of intentionally pushing or pulling the timing of a clock on a time axis as shown in FIG. 3 in order to use an internal clock adjusted to a valid window of data. At this time, the slave memory performs a setup/hold centering function of a data sampling clock by pushing or pulling the timing of a clock on the time axis in response to a setting signal.\nTypically, a delay line or delay chain is constructed with an inverter chain having a plurality of inverters as internal delay elements. The inverter chain is constructed with inverters connected in at least more than two levels having a relatively large amount of a unit delay. Therefore, it is not adequate to a case that requires a more precise delay.\nWhat is needed is a delay line or delay chain providing improved resolution degree by decreasing the amount of unit time delay and providing more precise control of the delay while minimizing skew in the clock signal.\nTherefore, the present invention is disclosed to solve the aforementioned problems and it is an object of the invention to provide an internal clock generating circuit to generate an internal clock precisely controlled as much as the necessary amount of delay in an improved resolution degree and a method for the same.\nIt is another object of the present invention to provide an internal clock generating method and the related circuit that can precisely sample data even when there is a skew between clock and data to be applied to a semiconductor memory.\nIt is a still another object of the present invention to provide an internal clock generating method and the related circuit that can control the delay time in response to an external signal.\nIt is a further another object of the present invention to provide an internal clock generating method and the related circuit that can minimize skew of an output clock, but in an improved resolution degree.\nIn order to accomplish the aforementioned objects in accordance with an aspect of the present invention, disclosed is an internal clock generating circuit of a semiconductor device comprising a delay chain having a plurality of delay units for generating multi-phase clocks by adjusting an input clock, a thermometer converter for outputting a thermometer code value in response to an input selection data, and a multiplexer for selectively outputting one of a plurality of clocks input from the delay chain in response to the thermometer code value.\nIn another aspect of the invention the multiplexer selectively outputs one of the plurality of clocks applied from the delay chain by dividedly multiplexing the thermometer code value into two stages of upper and lower bits.\nIn another aspect of the invention the delay chain additionally includes regenerators to restore a pulse form of a clock.\nIn another aspect of the invention the delay unit of the delay chain is constructed with RC delays.\nIn another aspect of the invention the regenerator is a short type pulse generator.\nIn another aspect of the invention a pulse regenerator is additionally included to restore a pulse form of a clock output from the multiplexer.\nIn another aspect of the invention the pulse regenerator is a short type pulse regenerator.\nIn another aspect of the invention the regenerators are respectively positioned in symmetry at positions where +/xe2x88x92 delay of the delay units against the thermometer code value changes non-linearly.\nIn another aspect of the invention the selection data is a binary code data applied from outside.\nIn another aspect of the invention the delay unit of the delay chain is constructed with path gates.\nDisclosed is a method of generating clock signals in a semiconductor device comprising the steps of generating multi-phase clock signals by adjusting an input clock signal through a delay chain having a plurality of delay units, decoding a thermometer code value in accordance with selection data, outputting one of the plurality of clock signals in response to the thermometer code value, and restoring a pulse shape of the output clock signal into its original state and outputting it as a delay-controlled internal clock signal.\nDisclosed is a semiconductor clock signal circuit, comprising means for generating multi-phase clock signals by adjusting an input clock signal through a delay chain having a plurality of delay units, means for decoding a thermometer code value in accordance with selection data, means for outputting one of the plurality of clock signals in response to the thermometer code value, and means for restoring a pulse shape of the output clock signal into its original state and outputting it as a delay-controlled internal clock signal."}
{"text": "In digital modulation schemes, data symbols are transmitted by modulating the amplitude and/or phase of a carrier wave having a certain frequency. For example, a data symbol typically represents an M-bit fragment of data, resulting in N=2M possible symbols. The set of N possible symbols are mapped to a set of N respective fixed complex numbers, which are referred to as constellation points and may be represented in the complex plane in the form of a constellation diagram. In order to transmit a given symbol, a complex carrier wave is multiplied by the value of the constellation point corresponding to the symbol, thereby modulating the amplitude and phase of the carrier by amounts corresponding respectively to the amplitude and phase of the constellation point.\nVarious constellations designs are used in various modulation schemes, including N-Quadrature Amplitude Modulation (QAM) in which the constellation comprises a square lattice of N regularly-spaced constellation points, and N-Phase Shift Keying (PSK) in which the constellation comprises a circular lattice of N regularly-spaced constellation points. Various other constellation designs are also known.\nIn order to measure the performance of a given constellation or between different constellations, various metrics may be used.\nFor example, capacity is a measure of the maximum rate of information that can be reliably transmitted over a communications channel. The maximum theoretical capacity of a channel is given by a well-known formula derived by Shannon. The Coded Modulation (CM) capacity is the maximum capacity achievable using a fixed non-uniform constellation without any coding constraints. The Bit Interleaved Coded Modulation (BICM) capacity is the maximum capacity achievable using a certain binary Forward Error Correction (FEC) scheme and fixed non-uniform constellation.\nIn addition, when comparing two systems, the difference in Signal-to-Nose Ratio (SNR) required achieving the same Bit Error Rate (BER) may be referred to as the SNR gain.\nIn contrast to uniform constellations, a non-uniform constellation is a constellation in which the constellation points are not regularly spaced. One advantage of using a non-uniform constellation is that performance may be increased, for example for SNR values below a certain value. For example, the BICM capacity may be increased by using a non-uniform constellation, when compared to an equivalent uniform constellation. Using a non-uniform constellation may also achieve a SNR gain over an equivalent uniform constellation.\nA constellation may be characterised by one or more parameters, for example specifying the spacing between constellation points. Since constellation points of a uniform constellation are regularly spaced, the number of parameters needed to characterise a uniform constellation is typically equal to 1. For example, for a QAM type constellation, the constellation is characterised by the (constant) lattice spacing. For a PSK type constellation, the constellation is characterised by the (constant) distance of each constellation point from the origin. On the other hand, since the spacing between constellation points in a non-uniform constellation varies, the number of parameters needed to characterise a non-uniform constellation is relatively high. The number of parameters increases as the order of the constellation (i.e. the number of constellation points) increases.\nOne problem with designing a non-uniform constellation is that a relatively high number of parameters need to be searched to find the optimum constellation. This problem is increased in the case of constellations of higher order. In the case of high-order constellations (e.g. constellations comprising more than 1024 constellation points), an exhaustive search across all parameters may be unfeasible.\nTherefore, what is desired is a technique for designing non-uniform constellations, and in particular, for designing non-uniform constellations for optimising performance (e.g. capacity and SNR performance). What is also desired is a technique for designing non uniform constellations using an algorithm having a relatively low complexity and relatively high computational efficiency."}
{"text": "The present invention relates to a yarn carrier, and a method and apparatus for manufacturing the same, and wherein the carrier is characterized by the ability to accommodate the tail of the yarn to be wound thereon between the carrier and the base plate of the winding machine, without severing the yarn tail during winding of the package.\nYarn carriers of the described type are designed to be wound with a yarn, and with a yarn tail extending from the carrier to permit the trailing end of the yarn on an exhausted carrier to be tied to the leading end of the yarn on a succeeding fully wound carrier. The yarn tail is typically formed at the beginning of the yarn winding operation by taking a length of yarn and extending it over the open large end of the carrier. The carrier is then mounted into a cradle of the winding machine which includes a base plate which fits into the large end of the cone and holds the yarn tail. A nose plate secures the small end of the carrier for proper rotation about a fixed axis in the machine, and the carrier is rotated by a rotating drum which engages the surface of the carrier and which feeds the yarn onto the rotating carrier in a predetermined reciprocating pattern.\nSuch yarn carriers are commonly manufactured from a sheet of paper which is wound about a mandrel to form a frusto-conical tubular member composed of several layers of the paper sheet. Both ends of the resulting tubular member are trimmed during the winding operation to provide even end surfaces, and the carrier is then finished to provide a rounded nose at the small end. The above end trimming operation is conventionally effected by a knife blade which moves radially inwardly against the paper sheet as it is being wound, and this operation inherently produces a rather sharp, annular burr at the intersection of the inner wall surface and the cut end surface at the large end of the tubular member. The annular burr at the large end is undesirable, in that it acts to sever the yarn tail when the yarn tail is positioned between the large end of the carrier and the base plate of a winding machine in accordance with the above described winding procedure. More particularly, the carrier often rotates relative to the base plate during the starting and stopping of the winding operation, and this relative movement causes the annular burr to sever the yarn tail.\nTo alleviate the severance problem, it has been proposed to polish the inside wall at the large end of the tubular member, to eliminate the burr. More particularly, this prior polishing operation has been performed with the use of a chuck having a profile matching the desired profile of the large end of the carrier, and it resulted in a slightly beveled edge on the inside of the large end of the carrier, and with the paper material at the large end being compressed to form a relatively hard inner surface.\nWhile the polishing of the inner end of the carrier has been a generally satisfactory method of removing the sharp annular burr and thereby avoiding the severance of the yarn tail during the winding process, modern winding machines have been designed with a universal base plate which is adapted to receive carriers of various angles of taper, to thereby avoid the expense of changing base plates whenever a different style of carrier is being wound. While such universal base plates are efficient from this point of view, they create a further problem in that the yarn tail is often pinched by a non-flush fit between the inside end surface of the carrier and the base plate, and such pinching in turn often results in the severance of the yarn tail. This problem is particularly acute where the carrier includes a hardened inner end surface as described above, since the hardened nature of the inner end surface tends to aggravate the severity of the pinching problem.\nOther solutions for the problem of severing the yarn tail have been proposed, note for example, U.S. Pat. Nos. 4,700,834 and 4,700,904, both issued to Martinez. These prior patents suggest the formation of spaced apart grooves, or grooves of crisscross configuration, or forming a ring of loose non-woven fibrous material of substantial thickness on the inside surface of the large end of the carrier. However, these rather elaborate constructions do not address the problem associated with the sharp annular burr as discussed above.\nIt is accordingly an object of the present invention to provide a yarn carrier, and a method and apparatus for manufacturing the same, which effectively avoids the above noted problem of severance of the yarn tail during the winding operation."}
{"text": "Portable handheld devices such as smartphones and tablet computing devices are widely distributed around the globe. In some markets, nearly every potential consumer carries their own device and uses it to send and transmit information. Devices are configured to use phone networks, like 3G and 4G networks, as well as data networks such as Internet connections accessed through Wi-Fi. Accordingly, these devices can offer their users the ability to access content almost constantly while they are turned on.\nThis ubiquitous and persistent connectivity of consumers offer obvious opportunities for gaming, marketing, news, communication, entertainment, and other electronic media. However, sophisticated and interactive electronic media such as interactive apps and games, group chats or “tweeting”, social networking apps with frequent updates, streaming music sites, or internet auction sites can quickly burden server capability as the volume of users increases. Further, old paradigms, such as creating a web page and posting it or giving people an app to play with on their phone, are being replaced by demand for continuous active content delivery such as games, apps, or media players that involve a person with real-time events, other parties, and new information on-the-fly.\nThe volume of content now involved in electronic media, the number of connected users, and the resulting demands on existing content delivery technologies leads to lag-times and increased latencies, as well as slow and sometimes stopped or broken data connections. Content preparation and delivery methods are poorly suited to the demands now being placed on them. As a result, users experience delays in obtaining access to content. Also, certain kinds of content requiring quick interactivity among high numbers of users is difficult to distribute at optimal rates."}
{"text": "The present invention relates to a tamper-resistant information processing device. It is particularly very effective when applied to cards such as the IC card.\nAn IC card is a device used for such purposes as to hold personal information which should not be altered without permission, to encrypt data by use of a cryptographic key (which is secret information), or to decrypt ciphertext. The IC card does not have any power source therein, but when it is inserted in a,reader/writer for IC cards, the IC card is supplied with power and becomes operable. When the IC card is in the operable state, it receives a command transmitted from the reader/writer, and carries out a process such as transfer of data according to the command.\nFIG. 1 shows a basic conceptual configuration of an IC card in which an IC card chip 102 is mounted on a card 101. As shown in the figure, an IC card generally has disposed thereon a supply voltage terminal Vcc, a ground terminal GND, a reset terminal RST, an input/output terminal I/O, and a clock terminal CLK. The positions of these terminals are specified in ISO International Standard 7816. The IC card receives power from the reader/writer and exchanges data with the reader/writer. Such communication between the IC card and the reader/writer is described, for example, on page 41 of a book entitled xe2x80x9cSMARTCARD HANDBOOKxe2x80x9d authored by W. Rankl and W. Effing and published by John Wiley and Sons in 1997.\nThe configuration of the semiconductor chip mounted on an IC card is basically the same as that of the ordinary microcomputer. FIG. 2 is a block diagram showing the basic configuration of the semiconductor chip mounted on an IC card. As shown in FIG. 2, the semiconductor chip for cards has a central processing unit (CPU) 201, a memory device 204, an input/output (I/O) port 207, and coprocessor 202. Some systems do not employ the coprocessor. The CPU 201 is a device for performing logic and arithmetic operations, while the memory device 204 stores programs and data. The input/output port is a device for communicating with the reader/writer. The coprocessor performs cryptographic processing itself or operations necessary for cryptographic processing at high speed. For example, types of coprocessors employed include a particular operation device for performing a residue operation for RSA and a cryptographic device for performing a rounding process for DES. There are many IC card processors which do not have any coprocessors. A data bus 203 is a bus connecting one device to another.\nThe memory device 204 includes such memories as a ROM (Read Only Memory), a RAM (Random Access Memory), and an EEPROM (Electric Erasable Programmable Read Only Memory). Information stored in a ROM cannot be altered, and therefore ROMs are used to store mainly programs. Information stored in a RAM, on the other hand, can be freely rewritten, but the stored information disappears once the power supply is interrupted. That is, since the power supply to an IC card is interrupted when the IC card is removed from the reader/writer, the RAM can no longer hold its contents after that. The EEPROM, in contrast, can continue holding its contents even if its power supply is interrupted. Therefore, the EEPROM is used for storing data which it is necessary to rewrite, and hold even when the IC card is removed from the reader/writer. For example, the number of the remaining call units of a prepaid telephone card is rewritten each time the card is used, and the call unit data must continue to be held even after the card is removed from the reader/writer. This is why the call unit data of the prepaid card is held in an EEPROM.\nThe present invention provides a tamper-resistant information device for use with cards having high security.\nSpecifically, an object of an embodiment according to the present invention is to reduce the correlation between the contents of data processing operations and consumed currents in a card component such as the IC card chip. Reducing the correlation between the contents of the data processing operations and the consumed currents in the chip makes it difficult to estimate what is being processed in the IC card chip and how, and to derive the cryptographic key from the observed waveforms of the consumed currents. Thus, the present invention provides cards with high security.\nSince IC cards have an IC card chip mounted thereon which is capable of holding programs and important information, they are used to store important information or internally perform cryptographic processing. It has been conventionally considered that the difficulty of breaking a code stored in an IC card is the same as the difficulty of deriving its encryption algorithm. However, it is pointed out that the details of the encryption processing operation and the cryptographic key may be derived by observing and analyzing the current consumed during the encryption process in the IC card, which may be easier than deriving of the encryption algorithm. The consumed current is obtained by measuring the current supplied from the reader/writer to the IC card. The details of this attack and its danger are described, for example, on page 263 (8.5.1.1 Passive Protective Mechanisms) of the book xe2x80x9cSMARTCARD HANDBOOKxe2x80x9d authored by W. Rankl and W. Effing and published by John Wiley and Sons. The following specifically describes the attack. Each CMOS constituting an IC card chip consumes a current when its output state switches from xe2x80x9c1xe2x80x9d to xe2x80x9c0xe2x80x9d or vice versa. Particularly, a large current flows through the data bus 203 when the bus value changes from 1 to 0 or vice versa. The current of the bus driver, the wiring employed, and the capacitance associated with transistors connected to the wiring cause such a current to flow. Therefore, it is possible to identify what is operating in the IC card chip by observing the consumed current.\nFIG. 3 shows single-cycle waveforms of currents consumed in an IC card chip. The current waveforms are different from one another as indicated by reference numerals 301 and 302, depending on the processed data. More specifically, such a difference occurs depending on data flowing through the bus 203 and data processed in the central processing unit 201.\nThe coprocessor 202 can perform, for example, 512-bit modular multiplication in parallel with the CPU processing. This means that it is possible to observe the waveform of a current different from that in the CPU for a long time. Therefore, the number of operations performed by the coprocessor can be measured by observing its particular current waveform. If the number of operations performed by the coprocessor has some relationship to the cryptographic key, it might be possible to derive the key from the number of the operations.\nFurther, if which operation is performed or what is operated by the coprocessor changes depending on the cryptographic key, the dependency might be found by observing the corresponding change in the consumed current, and the cryptographic key might be derived.\nSimilarly, in the CPU, the influence of each bit value of the cryptographic key on processed data might be obtained by changing the data a plurality of times and observing the corresponding change in each consumed current. It might be possible to derive the cryptographic key by statistically processing the waveforms of these consumed currents.\nThe ideas on which embodiments of the present invention are based include: dividing a process performed in an IC card so that attackers cannot specify the process as a whole; and inserting a dummy process. These methods make it difficult to identify the original process and derive the cryptographic key from the waveforms of the consumed currents.\nA tamper-resistant device as represented by the IC card chip is regarded as an information processing device having one or more data processing means which each comprise: a program storage unit for storing a program; a memory unit having a data storage unit for storing data; and a central processing unit (CPU) for performing a predetermined process to process data according to the program; wherein the program is composed of process instructions for giving an execution direction to the CPU.\nA method according to an embodiment of the present invention for scrambling the correlation between processed data and consumed currents in an IC card chip is to divide the data into pieces, and instead of performing a given operation(s) on the entire data as a whole, perform another different operation(s) on each piece of the divided data so as to still produce the same results as those that will be obtained if the given operation is performed on the entire data as a whole. As a result, the essential operation(s) can be concealed.\nSpecifically, pieces of scramble data R1, R2, . . . , and Rn are prepared. Original data D1 to be processed is divided into data blocks D1[1], D1[2], . . . , and D1[n].\nThe data blocks and scramble data are used to produce scrambled data blocks x[1], x[2], . . . , and x[n] by employing, for example, one of the following methods.\n(1) logical AND operation\n(2) x[1]=0, x[2]=xxe2x88x92v\n(3) x[1]=x AND R, x[2]=x AND xcx9cR, where xcx9cR is the inverse of R\nThat is, by using the scramble data R1, R2, . . . , and Rn, where R1 XOR R2 XOR . . . XOR Rn=2{circumflex over ( )}Lxe2x88x921 (L is the bit length of D1), the data block (original data) D1 is divided so that D2[1]=D1 AND R1, D2[2]=D1 AND R2, . . . , and D2[n]=D1 AND Rn, where n is an integer. In this case, the equation D2[1]+D2[2]+ . . . +D2[n]=D1 holds. In addition to the above logical AND operation, an ordinary addition operation or subtraction operation can be used for this purpose. A ring multiplication operation is performed on values obtained as a result of the above addition operation or subtraction operation to produce the final proper value. Since the randomly divided data blocks D2[1], D2[2], . . . , and D2[n] are used instead of directly using the original data D1, it is difficult to determine the original data D1 from information included in the observed current waveform alone. When a plurality of waveforms are statistically processed (for example, averaged to remove noise components from them), the characteristics of each waveform are eliminated, which further makes it difficult to determine original data (effectively hiding original information). It should be noted that the above randomly divided data may be produced through a division operation using pseudorandom numbers.\nAnother method for reducing the correlation between the contents of the data processing operation and the consumed current is to change the original data to be processed, and instead of performing a given operation on the original data, perform another different operation on the changed data so as to still produce the final proper results but consume a current different from that which will be consumed if the given operation is performed on the original data.\nSpecifically, random scalar data R for scrambling other data is prepared. Then, by using the prepared random scalar data R and a particular element V, the element to be processed is changed from x to x+R*V, where the symbols xe2x80x9c+xe2x80x9d and xe2x80x9c*xe2x80x9d denote ordinary addition and multiplication operations, respectively. The element V has the characteristic that whether or not the element V is added to data, an operation on the data produces the same results as if the element V were not added to the data. The above x+R*V can be used as an exponent or a scalar to scramble statistical processing of waveform observation data of consumed currents. It should be noted that the above element V acts as a number of 1 in a multiplication operation, and 0 in an addition operation. For example, when N=pq, where N is the modulus of a public key in the RSA cryptosystem, the element V is a multiple of (pxe2x88x921) (qxe2x88x921). When a scalar multiple of a base point on an elliptic curve is used, the element V is a multiple of the order of the base point.\nFurther, randomly determining the order in which each piece of the divided data is processed further makes it difficult to find the correlation between the contents of the data processing operation and the consumed current.\nStill further, combining all the above methods for scrambling encrypted data is effective in further reducing the correlation between the contents of the data processing operation and the consumed current.\nThe present invention can be applied to information hiding for modular multiplication and modular exponentiation in the RSA cryptography. Furthermore, in the elliptic curve cryptography, it can be applied to information hiding for multiplication and division in underlying fields, and the calculation of a scalar multiple of a base point. In modular multiplications, the logical AND operation described above is used to divide data, and then the distributive law is used to obtain the final proper result from the divided data. In the modular exponentiation and the calculation of a scalar multiple of a base point, the exponent is divided by means of ordinary subtraction, then modular exponentiation is performed on each piece of the divided exponent, and the product of the operation results is calculated to obtain the final result (proper answer). These operations are effective in scrambling encrypted data. It should be noted that the above modular multiplications include multiplication in a prime field."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an antenna assembly, and more particularly to an antenna assembly with a moveable antenna assembling on an electronic device, such as notebook.\n2. Description of the Prior Art\nSince the wireless communication technology of using electromagnetic wave to transmit signals has the effect of remote device transmission without cable connection, and further has the mobility advantage, therefore the technology is widely applied to various products, such as moveable phones, notebook computers, intellectual home appliance with wireless communication features. Because these devices use electromagnetic wave to transmit signals, the antenna used to receive electromagnetic wave also becomes a necessity in the application of the wireless communication technology. An antenna almost requires to receive and send signals in different directions. But the radiating performance of antennas inside the electronic devices is dissatisfactory due to the influence of components in the electronic devices. Outer antennas can eliminate the trouble, but the outer antennas can not achieve the handsome requirement of present designs of the electronic devices.\nHence, an improved antenna assembly with a moveable antenna is desired to overcome the above-mentioned shortcomings of the existing antennas."}
{"text": "When purchasing steel from multiple sources, quality of the steel may not be as dependable and as stable as when purchasing steel from a single source. Also, some steel has unusual composition such as alloying components like Boron, which reduces import duty but affects the behavior of the metal and requires different welding. When this chemical composition becomes hidden in the steel supply chain, the steel can be dangerous.\nOriginally, only chemical analysis could determine metallurgical components of the steel, and specialty steel suppliers employed metallurgists with laboratory equipment to test the steel. With the rapid availability of high performance microcomputers, heavy stand-alone machines use an arc for spectrographic analysis of a metal sample. Newer test equipment uses Energy Dispersive X-ray Fluorescence (ED-XRF) technology for spectral analysis of steel, although they are expensive and not practical for steel processing workshops.\nPlasma cutters have become common for metalworking, including working with steel, aluminum, and other metals. In plasma cutting, a plasma formed of a gas heated by an electric arc serves to conduct electricity into, and remove melted metal from, a metal workpiece. Plasma cutters may be used with “numerically controlled” (NC), computer controlled cutting machines, or may be handheld."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor package and a fabrication method thereof, and more particularly to semiconductor packages each with an extra plurality of electrical connecting points and a fabrication method thereof.\n2. Description of the Prior Art\nOwing to the trend toward multi-function, high-performance, and high-speed electronic products, semiconductor manufacturers nowadays are devoted to research and development of semiconductor devices integrated with multiple chips or packages with a view to meeting the requirement for today's electronic products.\nReferring to FIG. 1, U.S. Pat. No. 5,222,014 discloses a stack structure of a semiconductor package, and a method for fabricating the stack structure of a semiconductor package involves providing a first ball grid array (BGA) substrate 11 having bonding pads 110 disposed thereon, mounting a semiconductor chip 10 on the first BGA substrate 11, forming an encapsulant 13 on the first BGA substrate 11 such that the encapsulant 13 encapsulates the semiconductor chip 10, mounting and electrically connecting a second BGA substrate 12 (which has been packaged like the first BGA substrate 11) to the bonding pads 110 via solder balls 14.\nHowever, in the stack structure of the semiconductor package described above, the number of the bonding pads 110 electrically connecting the second BGA substrate 12 with the first BGA substrate 11 is restricted by the size of the encapsulant 13, thus limiting the type of semiconductor packages to be stacked and the number of I/O connections that can be formed, such that the type of semiconductor packages to be stacked and I/O connection layout on the second BGA substrate 12 would be restricted by the bonding pad 110 arrangement on the first BGA substrate 11. Moreover, due to the height limitation of the solder balls 14, the height of the encapsulant 13 disposed on the first BGA substrate 11 must be minimized (typically below 0.3 mm), which increases the difficulty of fabrication. Other stack structures of semiconductor packages such as that disclosed in U.S. Pat. Nos. 6,025,648 and 6,828,665 also experienced the same problem.\nBesides, in accordance with the foregoing known stack structures of semiconductor packages, electrical connecting points for forming electrical connections with external devices rely totally on the circuits on the substrate surface, while the encapsulant which occupies the majority of space in a package, however, is incapable of providing extra electrical connecting points, such that not only the overall electrical performance of the semiconductor product cannot be improved, the usage of the package would also be limited.\nThus, there is an urgent need to develop a semiconductor package and a method for fabricating the same, for providing an extra plurality of electrical connecting points, thereby solving the problem of package usage limitation, improving the electrical performance of electronic products and overcoming the stacking limitation in terms of size and type of semiconductor packages and the number of I/O connections."}
{"text": "For the purpose of conveying loose materials such as bulk materials, that is, rock/stones, mineral resources, excavation material, agricultural products, et cetera, use has long been made of troughed conveyor belts, which receive the conveyable material at a receiving location on their carrying side and discharge the same at a discharge location. Since the conveyable material is open to the environment as it is being transported, contaminants and environmental weathering influences can act on the conveyable material, and the latter can pollute the environment and also pose a risk to the environment. It is also the case, on account of their configuration, that troughed conveyor belts can be used to realize curves and gradients only to a limited extent. It is thus not usually possible, in conventional belt systems, to exceed an angle of inclination of 20° in gradient. If this is the limit of feasibility, it is necessary to connect a plurality of inter alia specific conveying belts with transfer locations. This increases the complexity for, and therefore the costs of, the conveyor system to a considerable extent.\nIn order to eliminate these disadvantages, conveyor belts which are closed during operation and are referred to as tube belts, tubular conveyor belts, pipe belts or mega pipes, were developed in the 1980s. The tube belts are rolled together between the receiving location and discharge location to give a closed tube, by virtue of the outer belt flanks overlapping and thus fully enclosing the conveyable material. This means that the conveyable material in the tube belt and the environment are completely separated from one another, since the tube belt remains closed over the conveying route. It is only for the purposes of receiving and discharging the transportable material that a tube belt widens and assumes the form of a conventional troughed conveyor belt. This rules out contamination of the bulk material along the conveying route and the associated environmental pollution. It is also the case that the conveyable material cannot be influenced by the environment during transportation. Further essential advantages of the tube belts in relation to the conventional troughed conveyor belts reside in the possibility of realizing very narrow three-dimensional curves and in the relatively high angles of inclination of up to 35° in gradient, which means that complicated three-dimensional curved routes can be realized by a single system. Since tube belts usually have a smooth surface on their carrying side, the angles of inclination are nevertheless limited to a gradient of up to 35°, depending on the bulk material properties.\nIn order to eliminate these disadvantages, conveyor belts which are closed during operation and are known as SICON®conveyor belts or pocket (conveyor) belts have also been developed. A pocket conveyor belt comprises two textile-reinforced profiles each with a steel cable vulcanized therein as a tension member. The profiles run over the sets of rollers and carry the pocket which accommodates the conveyable material. This droplet-shaped pocket consists of highly flexible rubber and is connected to the profiles by means of hot vulcanization. The profiles are arranged one above the other during transportation, and the belt is therefore closed off in a dust-tight manner. The belt is carried, and guided, by specific sets of rollers which, for the closed state of the belt, comprise a carrying roller and a guide roller. Further sets of rollers, each comprising a carrying roller and one to three guide rollers, are available for loading and unloading the belt and for curves and gradients.\nIn a manner similar to tube belts, the essential advantages of the pocket conveyor belt in relation to the conventional troughed conveyor belts reside in the possibilities of realizing very narrow three-dimensional curves and in the relatively high angles of inclination of up to 35°; in the case of conventional belt systems, the angle of inclination cannot usually exceed 20°. This makes it possible to realize complicated three-dimensional curved routes by a single system, without any transfer locations on the conveying route. In addition, the material in the pocket conveyor belt and the environment are completely separated from one another, since the pocket conveyor belt remains closed over the conveying route. For loading purposes, the pocket conveyor belt is opened with the aid of a specific set of rollers for opening and closing the belt. The belt can be unloaded at an overhead discharge point or an S-shaped discharge station. At the S-shaped discharge station, it is possible optionally for the belt to be emptied or for the conveyable material to be poured into the belt again.\nA pocket conveyor belt differs from the conventional tube belt not just in construction, but also in functioning and areas of application. It is thus possible for a pocket conveyor belt, depending on the profile size, to negotiate radii of 0.6 m or 1.0 m, which cannot be realized by a conventional tube belt. The minimum curve radius which can be realized by a tube belt is approximately 30 m. In contrast to the tube belt, the conveyor length, the conveyor cross section and the associated conveyor capacity and maximum possible material particle size of a pocket conveyor belt are very limited. All of this predestines a pocket conveyor belt for an “in-plant closed” transport of industrial bulk materials, while a tube belt is considered in practice to be more akin to an “out-plant closed” conveying principle for the entire range of particle sizes.\nFor a number of application cases, the advantages of the tube belts or the pocket conveyor belts and the steep conveyor belts are required at the same time, that is, a tube belt or pocket conveyor belt which can be used even at angles of inclination above 35°.\nU.S. Pat. No. 6,170,646, GB 1197700, U.S. Pat. No. 5,351,810, JP 480 48 385 U, JP 580 83 314 U, United States patent application publication 2012/0000751 A1, FR 14 968 97 A, GB 88 76 98 A, JP 582 16 803 A, U.S. Pat. No. 3,392,817 A and WO2005/085101 A1 disclose a number of technical solutions in this respect for increasing the angle of inclination of tube belts and pocket belts by differently shaped profiles having been applied to the carrying-side cover panel of a tube belt or pocket belt. The core idea of these approaches has been in each case, for elastic rubber or plastic-material strips connected to the conveyor belt to be fitted transversely to the longitudinal direction of the conveyor belt and to be offset at certain intervals from one another in the longitudinal direction. It is possible here for the transverse strips to span both the entire belt width and just part of the belt width. It can be established from these documents that the transverse strips may be configured both in a continuous state, in the form of ribs or wave-like strips, and in a divided state, for example, at right angles, in sawtooth form or in trapezoidal form. The divided transverse strips here are configured such that, when the tube belt or pocket belt is deformed in tube or pocket form, the flanks of the strip butt more or less against one another or overlap and thus form partition walls spaced apart in the longitudinal direction. Depending on the height, that is, radial formation, of the transverse strips, the conveyable material is retained in a force-fitting and form-fitting manner during transportation, and it is therefore possible to prevent the conveyable material from sliding back in the conveyor belt and thus to realize relatively large gradients. In the case of the transverse strips being virtually closed, it is even possible to realize vertical conveying directions, wherein purely form-fitting force transmission takes place.\nIt is a disadvantage of the above-described tube belts or pocket belts that they involve very high outlay, and are therefore expensive to produce. It is thus necessary for the transverse strips, on account of their size, in particular their radial extent, to be produced in the form of separate elements and to be applied subsequently to the conveyor belt for example by means of adhesive bonding, that is, by cold vulcanization. This requires the further operating steps of the transverse strips being separately produced and subsequently installed on the conveyor belt. Single-piece production of conveyor belts with transverse strips, that is, simultaneously with the vulcanization of the conveyor belt, is ruled out in production terms on account of the size of the transverse strips. It may also be necessary for the transverse strips to be installed on the conveyor belt for the first time at the site of use of the conveyor belt, so that there is no increase in the volume of the conveyor belt for transportation purposes. Furthermore, the adhesive-bonding locations constitute a weak point which, over time, will fail sooner than other constituent parts of the conveyor belt.\nIt is also disadvantageous that, if the known tube belts or pocket belts are suitable for relatively large angles of inclination, that is, above 35° in gradient, the conveyable material is retained in a form-fitting manner by the transverse strips and the latter are subjected to corresponding loading. This requires a corresponding stable and radial formation of the transverse strips with higher material and production costs than in the case of flatter profilings, although the latter do not allow such gradients. It is also the case that the higher transverse strips increase the transportation costs of the conveyor belts, because the latter cannot be wound as tightly for transportation purposes, that is, less belt length per rolled together belt drum can be transported in one journey. At the same time, this means that the pieces of belt which can be transported per drum in one journey are shorter, and there is therefore an increase in the outlay for installing the endlessly closed conveyor belts in the conveyor-belt system. As is also the case with conventional tube belts, the pocket conveyor belts have a smooth surface, as a result of which it is possible to realize the angles of inclination of up to 35°, depending on the bulk-material properties."}
{"text": "The present invention relates to control systems of the type having a plurality of remotely located process control units connected together through a communications link and, more particularly, to a control system in which one or more redundant control units each serve as a back-up for a plurality of the remotely located process control units.\nMany system type industrial installations, for example, those related to industrial process-type manufacturing and electrical power generation, employ a large number of physically distributed controlled-devices and associated sensors for effecting coordinated operation of the overall system. In the past, coordinated control of the various devices has been achieved by manual operation and various types of semi-automatic and automatic control systems including electromagnetic relay systems, hardwired solid-state logic systems, and various types of computer control systems. The computer systems have included central systems in which the various sensors and controlled devices are connected to a central computer; distributed control systems in which a remotely located computer is connected to each of the controlled devices and to one another, and hybrid combinations of the central and distributed systems. The successful functioning of the control system is vital to any industrial process, and, accordingly, distributed systems have generally been preferred over central systems because the failure of one of the remotely located control computers generally does not cause a system wide failure as in the case of the failure of the central computer in the central system. In copending application Ser. No. 115,161, filed Jan. 14, 1980, invented by Michael E. Cope and assigned to the assignee of this application, there is disclosed a distributed control system. The preferred embodiment of the present invention is employed in a distributed control system as disclosed in this copending application.\nIn the control system disclosed in application Ser. No. 115,161, U.S. Pat. No. 4,304,001, a plurality of remote process control units R.sub.n (remotes) are connected to various controlled devices and sensors and communicating with one another through a communications link, which transmits data serially. Each remote is assigned a unique succession number or position in a predetermined succession order with each remote unit assuming supervisory communication control of the communications link on a revolving or master for the moment basis in accordance with the remote's relative position in the succession order. Information transfer including process data and command control information is accomplished between a source remote R.sub.s and a destination remote R.sub.d by successively transmitting information blocks over the communications link with the destination remote R.sub.d testing the validity of the blocks and, if valid, responding with an acknowledgement signal (ACK), and, if invalid, a non-acknowledgement signal (NAK) is sent by the destination remote R.sub.d. The source remote R.sub.s will retransmit the information blocks in response to a non-acknowledgement signal from a destination remote.\nIn accordance with the present invention, one or more of the remotes connected to the communications link is a redundant remote. Each redundant remote is designed to monitor a plurality of other remotes which are referred to as primary remotes and each of which may perform active control operations. The redundant remote detects whether or not any one of the primary remotes has failed and if it has failed, it then will take over operation of the inputs and outputs of that remote by sending instructions or commands over the communications link to the failed remote ."}
{"text": "Data sources are currently scanned for a variety of purposes. For example, files can be scanned for viruses at predetermined locations. One example is the McAfee Vitran™ solution which combines heuristics and virus detection at predetermined locations and provides for isolation of suspect files. However, this solution is limited to virus scanning only.\nIn other solutions, specific data streams can be scanned for inappropriate content. For example, Microsoft's Internet Explorer™ provides for a Content Advisor that filters content based on user pre-selected criteria and rating placed on a web site. The Content Advisor in this case can filter content that creates fear, depicts drug or alcohol use, shows sexuality or nudity, among others. The filtering can be complete or limited based on the context of the web site. Other content scanning solutions include Net Nanny™ or Surfwatch™.\nThese solutions are, however, limited to one type of scanning and typically are performed on a specific data stream or file."}
{"text": "Tn U.S. Pat. No 3,612,793 issued on Oct. 12, 1971, to John O. Roeser and entitled \"Electrical Switch Components And Switches Formed Thereby,\" there is shown and described a unitary combination switchblade and contact means For compact switch mechanisms. As can be seen from FIGS. 8, 9 and 10 of the '793 patent, the combination switchblade and contact means 16 comprises a pair of identical members 109a (109b) each being formed of a first end portion 110 a mid-portion 116 and a curvilinear second end portion 118. The first end portion 110 consists of a pair of spaced apart legs 112 and 114. The legs are spaced apart by a dimension D so as to receive therebetween a tension spring 120, as illustrated in FIG. 5. A combination strut and attachment means 124 is shear formed in the curvilinear second end portion 118 and includes a reverse curve portion 131 for mating with the end of the tension spring 120 (FIG. 5).\nThe curvilinear second end portion 118 further includes an upper contact portion 128 and a lower contact portion 130 which is separated from the upper contact portion 128 by means of the reverse curve portion 131. The legs 112 and 114 are provided with sharpened edges 132 for pivoting within opposed V-shaped notch portions 106 formed within an elongated metal member 74 of the actuating means 114. The lower contact portion 130 is located within the same plane as the first portion 110 and the mid-portion 116 However, the upper contact portion 128 is off-set from that same plane. The lower and 15 upper contact portions 128 and 130 define movable contact members which are adapted to be snapped between opposed fixed contact portions upon upper terminal members 134 and lower terminal members 136, as illustrated in FIG. 5.\nThe upper portions of the upper terminal members 134 are formed with an off-set transverse portion 116 having a depending upper fixed contact 148 (FIGS. 14 and 15). The upper fixed contacts 148 serve to make electrical contact engagement with the upper contact portions 128 of the movable contact members. The lower terminal members 136 are provided with portions 156 which serve as lower fixed contacts for electrical contact engagement with lower contact portions 130 of the movable contact members.\nThe combination switchblade and contact means 16 is assembled with other components so as to form a compact precision, snap-action push-button switch 10 and is illustrated in FIGS. 2, 5 and 6. In operation, when the button 72 of the actuator means 14 is depressed so as to move downwardly, this causes the elongated metal member 74 to also more downwardly against the bias of the return springs 76 and 78. Initially, the combination 15 switchblade and contact means is restrained by means of the lower fixed contact portions 156 so that the lower contact portions 130 of the movable contact members remain in contact therewith.\nThis contact continues until the pivot points 106 pass the center line of the tension spring 120. At that time, the movable contact members will snap overcenter and the upper contact portions 130 thereof will become engaged with the upper fixed contacts 148. This engagement continues until such time when the pressure upon the top of the button 72 is released, thereby causing the return springs 76 and 78 to return the metal member 74 and button 72 to their rest positions shown in FIGS. 1 and 2. As the metal member 74 returns past the center line of the tension spring 120, the movable contact members will snap overcenter again in the reverse direction.\nThe present invention represents an improvement over the combination switchblade and contact means described above in connection with the '793 patent. The upper contact portion 128 and the lower contact portion 130 upon the second end portion 118 of the switchblade and contact means 16 in the aforenoted patent experience a high degree of stress during the switching operations which tends to shorten the service lives of such components. It would therefore be desirable to provide a novel and improved switch butterfly assembly for use in precision snap action switches which has a unique configuration so as to reduce the stresses upon the movable contact portions, thereby prolonging the service lives thereof."}
{"text": "Heretofore, reactive liquid polymer (RLP) epoxy adducts have been used for many years for toughening composites and adhesives. They provide outstanding improvement at room temperature, but are only minimally effective at lower temperatures. An epoxy adduct using CTB can improve low temperature performance, but the CTB, with no or low bound acrylonitrile content, is not miscible in uncured epoxy and will separate upon aging."}
{"text": "In the context of growing product functionalities of component carriers equipped with one or more electronic components and increasing miniaturization of such components as well as a rising number of components to be mounted on the component carriers such as printed circuit boards, increasingly more powerful array-like components or packages having several components are being employed, which have a plurality of contacts or connections, with ever smaller spacing between these contacts. In particular, a low loss electric connection of embedded components is desired. For example, component carriers with high I/O count and/or high-frequency signal transmission may be critical in this respect. Moreover, there is a strong tendency in the field of component carriers of continued miniaturization. In particular, it is desired to form electrically conductive connections with smaller and smaller dimensions. Thus, conventional procedures of contacting different layers of stacked layer-type component carriers reach their limits."}
{"text": "Endoscopes are increasingly being used in medical diagnosis and therapy When used as directed the endoscope becomes grossly soiled and massively infected with microorganisms which are present in body cavities, on the mucous membrane, and in the blood. Accordingly the instruments must be thoroughly cleaned and disinfected after each use. Endoscopes are precision instruments which are made from a combination of materials. They are difficult to clean in view of the sensitivity of the materials involved to chemical attack and because they have narrow lumens making access to and cleaning of interior surfaces difficult.\nUntil the last decade, it was common for soiled instruments to be placed on a towel or in a covered pan until they were sent to a centralised service where they were scrubbed, washed and either sterilised in a steam autoclave (if not heat sensitive) or chemically (e.g. with formaldehyde). In the last decade, there has arisen a particular concern to avoid transmission of very serious and sometimes fatal diseases such as may be carried in blood and tissue, for example hepatitis B, HIV, and other infections.\nNowadays, contaminated endoscopes and other medical instruments are typically treated in a first bath (“presoak” or “cleaning” bath) containing one or a combination of anionic and non ionic surfactants. The first bath may optionally include one, or a combination of enzymes, adapted to digest biological contaminants including cellular material, blood and other body fluids. Enzyme containing pre-soaking liquids are significantly more efficient in removing water insoluble and protein soils and are now considered the industry standard. In the case of surgical instruments requiring to be sterilised, the instruments are typically then removed from the first bath, washed free of enzyme solution and other residues, and then deposited in a second bath containing a chemical sterilizing agent (for example, glutaraldehyde). The first bath container is subsequently washed and then furnished with a fresh enzyme solution so that the process may be repeated. The necessity for separate cleaning and sterilizing baths arises since enzymes are denatured by all known sterilizing agents and since sterilizing agents are deactivated by enzymes (as enzymes are proteins). Accordingly it has to date proved impossible to provide a “single bath” cleaning and sterilizing treatment, although a two part system involving an enzyme treatment followed by addition of a phenolic disinfectant in the same bath has been proposed.\nSterilisation protocols are followed to prevent cross infection and therefore instruments used with one patient are not combined with those which may have been used with another in the presoak bath. It is noteworthy that the presoak is not passive. Staff are instructed to syringe detergent liquor through all the lumens, to brush biopsy channels, etc. A colonoscope requires up to 14 manual brushing-syringing-plugging-unplugging operations. It is usual for staff to wear latex gloves when handling instruments into or out of the baths and when performing such like operations.\nThe present inventor has observed that currently used procedures, while effective for preventing crossinfection between patients, in fact exposes medical and/or hospital staff to hitherto unrecognised health and safety risks. By virtue that the enzymes of the first bath digest the biological secretions holding the microorganisms, thus releasing them within the bath, and surfactants efficiently disperse them, the fluid content of the first bath is itself readily contaminated to high levels with infectious material. Contrary to the belief of some hospital staff, enzymes do not kill bacteria but rather release them. The present inventor has measured bacterial counts in excess of 1×109 colony forming units (“cfu”) per sq. cm. on instruments entering the first bath\nStaff are therefore at risk of infection (i) from splashes from the first bath either during scrubbing to release contaminants or during draining the first bath (or from splashes if an instrument is accidentally dropped into the bath), (ii) from glove failures (latex gloves have a “pinhole” failure rate of about 12%), (iii) from accidental glove immersion above the wrist line, (iv) from finger stick incidents in the bath resulting in glove and sometimes dermal penetration, (v) from aerosols created by brushes and syringes. In addition the wall surface of the first bath remains contaminated after the bath has been emptied and if not itself disinfected may be handled by unprotected staff. The last mentioned risk may be minimised by performing the sterilisation step immediately after the digestion step in the same container, but this does not avoid any of the other hitherto unrecognised risks and is wasteful in use of excess sterilant.\nIn some cases instruments may not be required to be sterilised, for example with spatulas, and holders which do not penetrate the body tissue, hair dressing implements and the like, it may be sufficient to disinfect the instruments to an appropriate standard. In such cases it would be desirable to provide a cleaning and disinfecting treatment capable of meeting the required standards with a single composition.\nAny discussion of the prior art herein is not to be construed as indicative of the state of the common general knowledge in the field.\nIt is the object of the invention to avoid or ameliorate the above discussed disadvantages of prior art, or at least provide a commercial alternative to the prior art.\nIt is an object of preferred embodiments of the present invention to avoid or at least ameliorate the risk of infection to persons cleaning medical instruments by such procedures.\nIt is a further object of at least some of the preferred embodiments to provide a single step cleaning and disinfecting composition for use in cleaning medical instruments.\nPreferred embodiments of the invention also address the risk of cross infection of instruments by virtue of multiplication of microorganisms, if any, which remain on the bath walls after each cycle of instrument cleaning.\nIt is an object of some embodiments of the invention to provide simple means for cleaning and disinfecting surfaces which require to be disinfected"}
{"text": "In electrostatographic imaging and recording processes such as electrophotographic printing, an electrostatic latent image is formed on a primary image-forming member such as a photoconductive surface and is developed with a thermoplastic toner powder to form a toner image. The toner image is thereafter transferred to a receiver member, e.g., a sheet of paper or plastic, and the toner image is subsequently fused or fixed to the receiver member in a fusing station using heat and/or pressure. The fusing station includes a heated fuser member which can be a roller, belt, or any surface having a suitable shape for fixing thermoplastic toner powder to the receiver member. Fusing typically involves passing the toned receiver member between a pair of engaged rollers that produce an area of pressure contact known as a fusing nip. In order to form the fusing nip, at least one of the rollers typically includes a compliant or conformable layer. Heat is transferred from a heated roller fuser member to the toner in the fusing nip, causing the toner to partially melt and attach to the receiver member.\nNormally included in a compliant heated fuser member roller is a resilient or elastically deformable base cushion layer (e.g., an elastomeric layer). The base cushion layer is usually covered by one or more concentric layers, including a protective outer layer. The base cushion layer is typically bonded to a core member included in the roller, with the roller having a smooth outer surface. Where the fuser member is in the form of a belt, e.g., a flexible endless belt that passes around the heated roller, it commonly has a smooth outer surface which may also be hardened. Similarly, a resilient base cushion layer can be incorporated into a deformable pressure roller used in conjunction with a relatively hard fuser roller.\nSimplex fusing stations attach toner to only one side of the receiver member at a time. In this type of station, the engaged roller that contacts the unfused toner is commonly known as the fuser roller and is a heated roller. The roller that contacts the other side of the receiver member is known as the pressure roller and is usually unheated. Either or both rollers can have a compliant layer on or near the surface. It is common for one of these rollers to be driven rotatably by an external source while the other roller is rotated frictionally by the nip engagement.\nIn a duplex fusing station, which is less common, two toner images are simultaneously attached, one to each side of a receiver passing through a fusing nip. In such a duplex fusing station there is no real distinction between fuser roller and pressure roller, both rollers performing similar functions, i.e., providing heat and pressure.\nIt is known that a compliant fuser roller, when used in conjunction with a harder or relatively non-deformable pressure roller, e.g., in a Digimaster 9110 machine made by Heidelberg Digital L.L.C., Rochester, N.Y., provides easy release of a receiver member from the fuser roller, because the distorted shape of the compliant surface in the nip tends to bend the receiver member towards the relatively non-deformable unheated pressure roller and away from the much more deformable fuser roller. On the other hand, when a conformable or compliant pressure roller is used to form the fusing nip against a hard fuser roller, such as in a DocuTech 135 machine made by Xerox Corporation, Rochester, N.Y., a mechanical device such as a blade is typically necessary as an aid for releasing the receiver member from the fuser roller.\nA conventional toner fuser roller includes a rigid cylindrical core member, typically metallic such as aluminum, coated with one or more synthetic layers usually formulated with polymeric materials made from elastomers. An elastically deformable or resilient base cushion layer, which may contain filler particles to improve mechanical strength and/or thermal conductivity, is typically formed on the surface of the core member, which core member may advantageously be coated with a primer to improve adhesion of the resilient layer. Roller cushion layers are commonly made of silicone rubbers or silicone polymers such as, for example, polydimethylsiloxane (PDMS) polymers disclosed, e.g., by the Chen, et al., patent (U.S. Pat. No. 6,224,978, assigned to Eastman Kodak Company, Rochester, N.Y.).\nThe most common type of fuser roller is internally heated, i.e., a source of heat is provided within the roller for fusing. Such a fuser roller generally has a hollow core member, inside of which is located a source of heat, usually a lamp. Surrounding the core member can be an elastomeric layer through which heat is conducted from the core member to the surface, and the elastomeric layer typically contains fillers for enhanced thermal conductivity [see for example the Fitzgerald patents (commonly assigned U.S. Pat. Nos. 5,292,606 and 5,336,539) and the Fitzgerald, et al., patent (commonly assigned U.S. Pat. No. 5,480,724)]. An internally heated fuser roller can be made using a condensation-polymerized silicone rubber material including solid filler particles, such as for example used in a NexPress 2100 digital color press (manufactured by NexPress Solutions LLC, Rochester, N.Y.).\nLess common is an externally heated fuser roller, which fuser roller is typically heated by surface contact with one or more heating rollers. An externally heated fuser roller can be made using an addition-polymerized silicone rubber material including solid filler particles. Externally heated fuser rollers are for example disclosed by the O'leary patent (U.S. Pat. No. 5,450,183, assigned to Eastman Kodak Company, Rochester, N.Y.), the Derimiggio, et al., patent (commonly assigned U.S. Pat. No. 4,984,027), the Aslam, et al., patent (commonly assigned U.S. Pat. No. 6,567,641), and the Chen, et al., patent (commonly assigned U.S. Pat. No. 6,490,430). Inclusion of thermal-conductivity-enhancing fillers enhances heat transfer from one or more external heating rollers typically used for the external heating of the fuser roller. Moreover, the thermal-conductivity-enhancing fillers enable intermittent use of an auxiliary heating device located within the roller.\nSome fuser rollers rely on film splitting of a low viscosity oil to enable release of the toner and (hence) receiver member from the fuser roller. The release oil is typically applied to the surface of the fuser from a donor roller coated with the oil provided from a supply sump. A donor roller is for example disclosed in the Chen, et al., patent (commonly assigned U.S. Pat. No. 6,190,771) which is hereby incorporated by reference.\nRelease oils (commonly referred to as fuser oils) are composed of, for example, polydimethylsiloxanes. When applied to the fuser roller surface to prevent the toner from adhering to the roller, fuser oils may, upon repeated use, interact with PDMS material included in the resilient layer(s) in the fuser roller, which in time can cause swelling, softening, and degradation of the roller. To prevent these deleterious effects caused by release oil, a thin barrier layer made of, for example, a cured fluoroelastomer and/or a silicone elastomer, is typically formed around the resilient cushion layer, as disclosed in the Davis, et al., patent (U.S. Pat. No. 6,225,409, assigned to Eastman Kodak Company, Rochester, N.Y.) and the Chen, et al., patents (U.S. Pat. No. 5,464,698, and 5,595,823, assigned to Heidelberg Digital L.L.C., Rochester, N.Y.). A fluoro-thermoplastic random copolymer outermost coating can also be used for this purpose, as disclosed in the Chen, et al., patents (commonly assigned U.S. Pat. Nos. 6,355,352 B1 and 6,361,829 B1).\nTo rival the photographic quality produced using silver halide technology, it is desirable that electrostatographic multicolor toner images have high gloss. To this end, it is desirable to provide a very smooth fusing member contacting the toner particles in the fusing station. A gloss control outer layer (which also serves as a barrier layer for fuser oil) can be provided as disclosed in the Chen, et al., patent application (commonly assigned U.S. patent application Ser. No. 09/608,290). A fluorocarbon thermoplastic random copolymer useful for making a gloss control coating on a fuser roller is disclosed in the Chen, et al., patent (commonly assigned U.S. Pat. No. 6,429,249).\nIn the fusing of the toner image to the receiver member, the area of contact of a conformable fuser roller with the toner-bearing surface of a receiver member sheet as it passes through the fusing nip is determined by the amount of pressure exerted by the pressure roller and by the characteristics of the resilient cushion layer. The extent of the contact area helps establish the length of time that any given portion of the toner image will be in contact with and heated by the fuser roller. It is generally advantageous to increase the contact time by increasing the contact area so as to result in a more efficient fusing process. However, unless the effective modulus for deforming a compliant roller in the nip is sufficiently low, high nip pressures are required to obtain a large nip area. Such high pressures can be disadvantageous and cause damage to a deformable roller, e.g., such as pressure set or other damage caused by edges of thick and/or hard receiver members as they enter or leave the nip. Hence a low modulus deformable roller is desirable.\nIt is known from the Chen, et al., patent (commonly assigned U.S. Pat. No. 5,716,714) that use of a relatively soft deformable fusing-station roller (e.g., a deformable pressure roller having a low effective modulus for deformation) can advantageously reduce the propensity of a fusing station nip to cause wrinkling of receiver members passing through the nip.\nOne way to try to create a low modulus fusing-station roller is to use a foamed material, e.g., a cured material having an open-cell or a closed-cell foam structure, with the material inclusive of suitable strength-enhancing and/or thermal-conductivity-enhancing fillers. Attempts to utilize such foamed materials, for example as base cushion layers, have not generally been successful, for a number of reasons. The thermal conductivity of closed-cell structures tends to be disadvantageously low, even when squeezed in a fusing nip. Although an open-cell structure can be squeezed relatively flat in a fusing nip, the resilience typically becomes compromised because opposite walls within the foam tend to stick together under the heat and pressure of the nip. Moreover, foamed polymeric materials generally have poor tear strength, and shear strength also tends to be low. As a result, fusing-station rollers incorporating foamed base cushion layers are quite susceptible to damage and tend to age rapidly.\nIn particular, attempts to make fusing-station rollers with fluoroelastomer foamed materials, which have desirable low surface energy and high thermal stability, have not been successful because of the tendency to incur a “pressure set” under the high loading typically present in fusing station nips. For example, foam rollers made with VITON® fluoroelastomers are susceptible to “pressure set”.\nSuitable thermal conductivity of synthetic layers used in fusing-station rollers is attainable by the use of one or more suitable particulate fillers, the thermal conductivity being determined by the filler concentration. The thermal conductivity of most polymers is very low and the thermal conductivity generally increases as the concentration of thermally conductive filler particles is increased. However, if the filler concentration is too high, the mechanical properties of a polymer are usually compromised. For example, the stiffness of the synthetic layers may be increased by too much filler, e.g., so that there is insufficient compliance to create a wide enough nip for proper fusing. Moreover, too much filler will cause the synthetic layers to have a propensity to delaminate or crack or otherwise cause failure of the roller. Because the mechanical requirements of fusing-station rollers require that the filler concentrations generally be moderate, the abilities of the layers to transport heat are thereby limited. In fact, the total concentration of strength-enhancing and thermal-conductivity-enhancing in prior art internally heated compliant fuser rollers has reached a practical maximum. As a result, the number of copies that can be fused per minute is limited, and this in turn can be the limiting factor in determining the maximum throughput rate achievable in an electrostatographic printer.\nAn auxiliary internal source of heat may optionally be used with an externally heated fuser roller, e.g., as disclosed in the Aslam, et al., patent (commonly assigned U.S. Pat. No. 6,567,641) and in the Chen, et al., patent (commonly assigned U.S. Pat. No. 6,490,430). Such an internal source of heat is known to be useful when the fusing station is quiescent and/or during startup when relatively cold toned receiver members first arrive at the fusing station for fusing therein. In order for such an auxiliary internal source of heat to be effective (when intermittently needed) the fuser roller must have a sufficiently large thermal conductivity. However, this requirement conflicts with a need to keep heat at the surface of an externally heated fuser roller, i.e., so as not to unnecessarily conduct heat into the interior which would compromise the fusing efficiency of the roller. On the other hand, it is important to have a high enough thermal conductivity at the surface of the fuser roller to ensure efficient transfer of heat to the fuser roller from one or more heating rollers contacting the surface. Moreover, in order to have high efficiency, externally heated fuser rollers rely to a certain extent on thermal conduction of heat around the surface of the roller.\nWays to improve upon the above-described limitations associated with externally heated elastically deformable fuser rollers (including an optional auxiliary internal source of heat) are disclosed in the Chen, et al., patent applications (commonly assigned U.S. patent application Ser. Nos. 10/139,486 and 10/139,464). In the Chen, et al., U.S. patent application Ser. No. 10/139,486, an externally heated fuser roller having improved efficiency includes a core member, a base cushion layer around the core member, a relatively thin heat storage layer around the base cushion layer, and an outer gloss control layer around the heat storage layer, wherein the heat storage layer is loaded with more thermally conductive filler than is the base cushion layer and hence has a higher thermal conductivity. In the Chen, et al., U.S. patent application Ser. No. 10/139,464, a thin heat distribution layer is further included between the heat storage layer and the outer gloss control layer. While the fusing efficiencies relating to U.S. patent application Ser. Nos. 10/139,486 and 10/139,464 are much improved, the fuser rollers (respectively having three-layer and four-layer structures around the core member) are relatively expensive to manufacture, and may be susceptible to delamination with prolonged use.\nIt is known that instead of solid fillers, certain hollow fillers can be included in an addition-polymerized silicone rubber for the purpose of lowering rather than increasing the thermal conductivity of a deformable fuser roller, as disclosed in the Meguriya patent (U.S. Pat. No. 6,261,214, assigned to Shin-Etsu Chemical Company, Ltd., Tokyo, Japan). In particular, the Meguriya patent discloses incorporation into the silicone rubber of hollow filler particles (also known as microballoons) manufactured under the tradename EXPANCEL® available from Expancel, Sundsvall, Sweden, and Duluth, Ga.\nHollow microballoons are well known and are disclosed for example in the Morehouse, et al., patent (U.S. Pat. No. 3,615,972, assigned to Dow Chemical Company, Midland, Mich.). Microballoons are made from thermoplastic microspheres which encapsulate a liquid blowing agent, typically a hydrocarbon liquid. Such microspheres are made in unexpanded form. The walls of the unexpanded microspheres are generally impermeable to the liquid blowing agent, i.e., diffusion of molecules of the liquid blowing agent through the walls is typically negligible. An expanded form of a microsphere, i.e., a microballoon, is obtained by heating an unexpanded microsphere to a suitable temperature so as to vaporize the blowing agent, thereby causing the microsphere to grow to a much larger size. Too high of a heating temperature can result in some loss of internal vapor pressure and a shrinking of the microballoon. Methods for expanding microspheres are disclosed in numerous patents, such as, for example, the Gunderman, et al., patent (U.S. Pat. No. 3,914,360, assigned to Dow Chemical Company, Midland, Mich.), the Edgren, et al., patent (U.S. Pat. No. 4,513,106, assigned to KemaNord AB, Stockholm, Sweden) and the Morales, et al., patent (U.S. Pat. No. 6,235,801 B1). Expansion is generally irreversible after cooling, and the expanded microballoon form is stable under normal ambient conditions and can be sold as a dry powder or alternatively as a slurry in a liquid vehicle. Expanded microspheres or microballoons which are available commercially can be incorporated into various materials, such as for example to make improved paints or lightweight parts. Unexpanded microspheres are also available commercially for incorporation into various types of materials (e.g., expandable inks) or for manufacture of solid parts, e.g., by thermal curing in a mold so as to expand the microspheres. The shell material of certain microsphere particles can include finely divided inorganic particles, e.g., silica particles.\nThe use of microspheres in a compressible layer of a digital printing blanket carcass is disclosed in the Castelli, et al., patent (U.S. Pat. No. 5,754,931, assigned to Reeves Brothers, Incorporated, Spartanburg, S.C.). The microspheres are uniformly distributed in a matrix material which includes thermoplastic or thermosetting resins.\nThe Dauber, et al., patent (U.S. Pat. No. 5,916,671, assigned to W.L. Gore & Associates, Incorporated, Newark, Del.) discloses a resilient gasket made of a composite of polytetrafluoroethylene and resilient expandable microspheres.\nThere remains a need to provide for an electrostatographic machine an improved fusing station having high productivity which includes fusing-station members that are simple in construction, are very durable, and are made of material that can resist gouges or pressure damage from edges of receiver members moving through a high pressure fusing nip.\nSpecifically, there remains a need for a tough, long lasting fuser roller which can have only one layer coated on a core member, and which is thereby simple in structure by comparison with multi-layer prior art fuser rollers. This layer is required to be chemically unreactive, stable at high temperatures, and impervious to fuser oil. Moreover, there remains a need for an improved pressure roller having a similarly simple structure and which has similar properties.\nA crosslinked fluoroelastomer is a desirable material for making fuser rollers and pressure rollers, because of low surface energy, chemical inertness, imperviousness to fuser oil, and high-temperature stability."}
{"text": "Digital storage oscilloscopes are well known and have many capabilities for displaying waveforms of interest to a user and information related to those waveforms. For example, the envelope of a waveform is sometimes of interest to a user. Oscilloscopes exist which have an operating mode, termed `envelope mode`, in which the minimum and maximum excursions of a signal (termed the envelope of the signal) over a predetermined number of preceding signal acquisition periods is displayed. At each acquisition, the acquired signal is analyzed to determine the new envelope waveform, and that waveform is displayed on the screen. During display of the envelope waveform, the oscilloscope is unable to display the waveform of the underlying signal. In such an oscilloscope, the envelope waveform is occasionally reset, and the signal being monitored is briefly displayed to provide the user an idea of the underlying signal. Such oscilloscopes display the envelope waveform as a fully filled-in waveform. Such a display, however, provides no information concerning the underlying signal.\nFor example, FIG. 1 is a display screen 2 of a prior art oscilloscope screen displaying the envelope waveform 4 of an underlying signal (not shown). In FIG. 1, the background 5 of the display screen is dark, and the displayed signal 4 is represented by a white display, in the known manner. The signal being monitored has the envelope 4 illustrated in white on in FIG. 1. As can be seen, the fully filled-in displayed waveform 4 gives no indication of the underlying signal, but only displays the minimum and maximum excursions of that signal over the period of time over which the underlying signal is enveloped.\nOscilloscopes also have the ability to rasterize an acquired waveform using sparsely populated vectors with pixels displayed at a reduced display intensity. For example, see U.S. patent application Ser. No. 09/026 185, filed &lt;&lt;filing date&gt;&gt; by Sullivan et al., incorporated by reference herein. By rasterizing using sparse vectors, and with reduced intensity, the rasterizer can process more acquisitions, giving a more accurate representation of the signal being observed.\nThough the envelope mode is useful, a user may wish simultaneously to view the underlying signal and its envelope. An oscilloscope which can display the envelope waveform, while simultaneously displaying sufficient acquisitions of the underlying signal to provide an accurate indication of that signal is, thus, desirable."}
{"text": "SAR (specific absorption rate) for users of portable wireless devices (PWDs) is a matter of increasing concern. RF radiation to the user's head results from the free-space generally omnidirectional radiation pattern of typical current PWD antennae. When PWDs equipped with such an antenna are placed near the user's head, the antenna radiation pattern is no longer omnidirectional as radiation in a large segment of the azimuth around the user is blocked by the absorption/reflection of the user's head and hand. An antenna system for PWDs that greatly reduces radiation to the body and redirects it in a useful direction is also desirable.\nPrior art antennas for PWDs may cause audio noise in a hearing aid of the user. Referring to FIG. 16, a diagrammatic view of a prior art PWD 400 (in the form of a cellphone) used in the vicinity of a hearing aid 402 is illustrated. Cellphone 400 has a speaker on the keyboard surface near the top of the phone, which is normally aligned with the center of the user's ear 404 during use. Hearing aid 402 may be any type, including in-ear and behind-ear variations. Hearing aid 402 has an amplified audio output port 406, which is inserted into the ear canal of the ear 404. During operation, an electromagnetic field 408 is generated around cellphone 400 by omnidirectional antenna 440. In operation, electromagnetic field 408 illuminates the hearing aid 402, user's ear 404, and the user's head. RF noise is induced in the hearing aid by the field 408, resulting in excessive audio noise being presented to the user.\nThe planar inverted F antenna or PIFA is characterized by many distinguishing properties such as relative lightweight, ease of adaptation and integration into the device chassis, moderate range of bandwidth, omni directional radiation patterns in orthogonal principal planes for vertical polarization, versatility for optimization, and multiple potential approaches for size reduction. Its sensitivity to both vertical and horizontal polarization is of practical importance in mobile cellular/RF data communication applications because of the absence of the fixed antenna orientation as well as the multi-path propagation conditions.\nTo assist in the understanding of a conventional PIFA, a conventional single band PIFA assembly is illustrated in FIG. 17. FIG. 17 illustrates a prior art single-band PIFA antenna 440 located on the rear side 442 of a personal wireless device 444. PIFA 440 consists of a radiating element 446, a ground plane 448, a feed conductor 450, and a grounding conductor 452. PIFA 440 is typically positioned near an upper edge of ground plane 448 with the free end of radiating element 446 being closer to a user's hand than the feed conductor 450 and grounding conductor 452. The feed conductor 450 serves as a feed path for radio frequency (RF) power to the radiating element 446. The feed conductor 450 is electrically insulated from the ground plane 448. The grounding conductor 452 serves as a short circuit between the radiating element 446 and the ground plane 448. The resonant frequency of the PIFA 440 is determined by the length (L) and width (W) of the radiating element 446 and is slightly affected by the locations of the feed conductor 450 and the grounding conductor 452. The impedance match of the PIFA 440 is achieved by adjusting the dimensions of the conductors 450, 452, and by adjusting the separation distance between the conductors 450, 452. In operation, ground plane 448 radiates RF energy which is absorbed by a user's hand. Antenna 440 can be configured to reduce the SAR value to 1.6 mw/g with the PWD 444 transmitting at the 0.5 watt cw level. However, even at this level audio noise may be generated in a user's hearing aid by operation of PWD 444. Another limitation of the PIFA is its relatively low front-to-back ratio. Front-to-back ratios of typically PIFAs range from 0 to 2 dB. A 5 dB front-to-back ratio may be achieved by substantially increasing the distance between radiating element 446 and ground plane 448. A need exists for an antenna exhibiting substantially greater front-to-back ratios.\nFIG. 18 illustrates a prior art dual-band PIFA antenna 462, which is located on the rear of a personal wireless device 464, and electrically connected to ground plane 466 at one end and capacitively coupled to ground plane 466 at another end. PWD 464 further includes a battery pack 470 positioned away from antenna 462. In normal operation, PWD 464 is oriented in an upright manner so that end 472 is generally above end 474. Ground plane 466 is provided by the ground traces of the printed wiring board (PWB). The portion of antenna 462 indicated by numeral 476 resonates over a higher frequency band, while the entire portion 476, 478 of antenna 462 resonates over a lower frequency band. PIFA antenna 462 is grounded at its upper end at location indicated as numeral 480 to ground plane 466. PIFA antenna 462 is capacitively coupled at pad 482 in a direction away from upper end 472 of PWD. This type of antenna provides some reduction in SAR, but has limited ability to reduce hearing aid noise from a digital PWD.\nDespite all of the desirable properties of a PIFA, the PIFA has the limitation of a rather large physical size for practical application. A conventional PIFA should have the semi-perimeter (sum of the length and the width) of its radiating element equal to one-quarter of a wavelength at the desired frequency. With the rapidly advancing size miniaturization of the radio communication devices, the space requirement of a conventional PIFA is a severe limitation for its practical utility."}
{"text": "This invention relates to pens which utilize a fluid ink and whose writing tips are axially positionable with respect to an outer sheath.\nVarious types of pens using an aqueous or non-aqueous fluid are in widespread use for writing, drawing, painting or marking purposes, and may collectively be referred to as marking pens. Said fluids, which may be inks having soluble dyes, or paints having dispersed pigments, may be generically referred to as marking fluids. Many of such pens are provided with protective caps which prevent evaporation of the fluid and prevent accidental contact of the fluid with the clothing or skin of the user or with other objects. However, the placement and removal of the cap is troublesome, and the cap is frequently misplaced.\nMarking pens which avoid the need for a protective cap are well known and generally employ a mechanism whereby the tip of the pen can be retracted into a protective enclosure within an elongated sheath comprising the outer body of the pen. Marking pens of such construction are disclosed for example in U.S. Pat. Nos. 4,218,154; 3,652,172; and 4,540,300. The protective enclosures and associated retracting mechanisms are, however, generally of complex, expensive construction and do not endure long term use.\nAccordingly, it is an object of the present invention to provide a retractable capless marking pen having a protective enclosure within the body of the pen that prevents evaporation of marking fluid from the tip of the pen.\nIt is a further object of this invention to provide a marking pen as in the foregoing object of rugged and durable construction amenable to low cost manufacture.\nThese objects and other objects and advantages of the invention will be apparent from the following description."}
{"text": "U.S. patent application No. 528,043, filed Aug. 31, 1983, teaches the use of a flexible magnetized layer (e.g., the material commercially designed \"Plastiform\" available from Minnesota Mining and Manufacturing Company, Saint Paul, Minn.) to magnetically hold and drive flexible sheets of coated abrasive material containing ferromagnetic material. Typically the magnetized layer is incorporated in a back up pad wherein it is attached to one surface of a layer of resiliently compressible foam, the back up pad either being adapted to be manually manipulated or having a rigid backing plate adapted to be driven by a drive motor attached to an opposite surface of the layer of foam so that the motor can be used to drive the pad and thereby the abrasive against a surface to be abraded, while the layer of foam affords applying resilient pressure between the abrasive and the surface to be abraded.\nThe flexible magnetized layer can securely hold the abrasive sheets in place to abrade a workpiece while the back up pad is being used manually or driven by the motor, however, the material of the magnetized layer is somewhat friable and thus the edge of the magnetized layer can be worn away during use."}
{"text": "In the related art, there is known an electronic apparatus (such as a mobile phone terminal) where a card unit (such as a SIM card) obtained by installing various devices on a substrate can be mounted. Such an electronic apparatus is driven by electric power from a battery. In addition, the electric power from the battery is also supplied to the card unit mounted on the electronic apparatus.\nHowever, there is a problem in that it is difficult to drive the devices installed in the card unit when an electric power supply capability of the battery becomes lower than a predetermined value (for example, when a remaining battery level is lower than a predetermined value and the like)."}
{"text": "1. Field of the Invention\nThe present invention relates to an application providing method for a mobile terminal and, more particularly, to a mobile terminal and an application providing method that provide an application using an application package installer having multiple pieces of signature information.\n2. Description of the Related Art\nWith active development of applications for mobile terminals, application servers provide users with various application packages in downloadable form.\nTo install an application onto a mobile terminal using an application package, the mobile terminal has to receive signature information signed with a signature key suitable for the mobile terminal together with the application package.\nCurrently, an application package is associated with a single piece of signature information.\nTo provide an application to a particular mobile terminal, multiple application packages are created corresponding in number to terminal types and each such application package is associated with signature information specific to a particular terminal type.\nCreating multiple application packages for different terminal types is an unnecessary use of time and effort. In addition, whenever a new type of mobile terminal is manufactured, new application packages need to be created and maintained."}
{"text": "Field of the Disclosure\nEmbodiments described herein generally relate to a chamber liner for a semiconductor process chamber and a semiconductor process chamber having a chamber liner. More specifically, embodiments disclosed herein relate to a chamber liner for processing temperatures greater than about 650 degrees Celsius while shielding chamber components from halogen damage.\nDescription of the Related Art\nIn the fabrication of integrated circuits, deposition processes such as chemical vapor deposition (CVD) or plasma enhanced CVD processes are used to deposit films of various materials upon semiconductor substrates. These depositions may take place in an enclosed process chamber. The process gases used in the depositions deposit films on the substrate, but may also deposit residue on the internal walls and other components of the process chamber. This residue builds up as more substrates are processed in the chamber and leads to generation of particles and other contaminants. These particles and contaminants can lead to the degradation of the deposited films on the substrates causing product quality issues. Process chambers must be periodically cleaned to remove the deposited residue on the chamber components.\nA chamber liner may be disposed in the chamber to define a processing region in a desired location within the chamber with respect to the substrate. The chamber liner may be configured to assist in confining the plasma to the processing region and help prevent other components in the chamber from being contaminated with deposited materials, such as the residue mentioned above. The process gases may be supplied above a substrate support. A purge gas may be provided from below the substrate support to prevent process gases from reaching areas at the bottom of the chamber and causing deposit of residue in the areas below the substrate support. The process gas and the purge gas may be removed from the process chamber using a common exhaust disposed away from the process area, such as around an outer perimeter of the process chamber, to prevent mixing of the purge gas with the process gas in the process area. Using the arrangement described above, particle formation can occur in the process area above the substrate and cause defects in the products made in the process chamber.\nFurthermore, substrate processing temperatures are typically capped between about 400 degrees Celsius and about 480 degrees Celsius for silicon based depositions due to the aggressive erosion and corrosion by the halogen clean on the high temperature components. As such, optimal film quality is often sacrificed due to manufacturability and reliability concerns.\nThus, there is a need for an improved liner for a process chamber to prevent particle formation and/or to permit significantly higher substrate processing temperatures while shielding sensitive components from halogen damage."}
{"text": "1. Technical Field\nThe present invention is related to a vehicle mirror apparatus in which the mirror-face angle of a mirror is adjusted by a worm gear being rotated to rotate a wheel gear.\n2. Related Art\nIn a mirror-face adjusting actuator described in Japanese Patent Application Laid-Open (JP-A) No. 2006-88788, a worm gear is connected to the rotation shaft of a motor, the worm gear is meshed with a worm wheel, and the worm wheel is rotated by rotating the worm gear using motor driving force, thereby tilting a side mirror. The mirror-face angle of the side mirror is accordingly adjusted.\nHowever, with the mirror-face adjusting actuator, and in particular in circumstances in which tilting of the side mirror is stopped and rotation of the worm wheel is in a stopped state, the worm gear greatly displaces in the radial direction towards the worm wheel radial direction outside when rotation force due to driving force of the motor acts on the worm gear. This leads to the possibility of operation noise of the mirror-face adjusting actuator getting louder when the meshing depth of the worm gear with the worm wheel changes greatly."}
{"text": "The present invention relates to the removal of sediment from bodies of water, and in particular to a sediment removal method and system that removes sediment from significant depths while limiting or avoiding turbidity.\nMethods of removing sediment from bodies of water exist, however, such known methods of removal often include undesirable side effects. For example the equipment and process of dredging generates significant turbidity in that the sediment is openly disturbed from its settled condition. Such disruption to the sediment bed causes portions of the sediment to become suspended in the water. This is undesirable, particularly when the sediment is contaminated, as the level of contamination in the water is increased. Furthermore, dredging processes, such as cutter head dredging, are inefficient in that a much greater percentage of water is removed as compared to sediment. Cutter head dredging is also an unfocused or less controlled process in that the equipment is not easily maneuvered and, therefore, also susceptible to damage from submerged objects.\nTherefore an apparatus is needed that enables controlled removal of sediment while avoiding turbidity."}
{"text": "The Superconducting Quantum Interference Device, or SQUID, is well known as a sensitive detector of weak magnetic fields. As indicated in FIG. 1, a SQUID is comprised of a superconducting loop containing one or more Josephson junctions (indicated by X in FIG. 1), and magnetic flux Φ is inductively coupled into the loop through a coupling inductor L. A Josephson junction is known to act as a lossless nonlinear inductance below its critical current Ic, and also exhibits nonlinear resistance above Ic. A one-junction SQUID (FIG. 1A) comprises a single junction and a loop, and exhibits a nonlinear impedance which depends on the flux Φ in a periodic manner, with periodicity Φ0=h/2e=2.07 fT−m2, where h is Planck's constant and e is the charge on the electron. The one-junction SQUID does not have a direct voltage readout since the junction is shunted by a lossless superconducting inductor, so it must be embedded in a radio-frequency (RF) circuit for its impedance to be measured. For this reason, this structure is sometimes called an RF-SQUID, although the flux Φ can be at any frequency down to DC. Another SQUID device is the two-junction SQUID (FIG. 1B), which is generally operated with a DC bias current greater than the device critical current Ic=Ic1+Ic2 of the two constituent junctions. This then exhibits a DC voltage output across the Josephson junctions, modulated by the flux Φ in a way that is again periodic in Φ0 (FIG. 1C). This two-junction SQUID was historically called the DC-SQUID, since it operates down to DC frequencies, although it may alternatively operate with flux modulation up to gigahertz radio frequencies. The DC SQUID is much more commonly used than the RF SQUID, so in general usage, the term SQUID commonly refers to two-junction SQUID. Such a SQUID may be used not only for low-frequency magnetic field detectors, but also for radio-frequency amplifiers and active radio antennas if appropriate inductive inputs are used.\nBoth high-Tc and low-Tc superconductor based SQUID-amplifiers have been studied during the past ten years [7], [8], [9], [15], [16], [17]. See also Hilbert, U.S. Pat. No. 4,585,999, “RF amplifier based on a DC SQUID”. However, the characteristics of the amplifiers are still far from desired performance values. Despite the fact that the noise temperature Tn≈1-3 K [15] is reasonably low, the dynamic range (amplitude ratio) D=(Tsat/Tn)1/2 of the amplifiers is strongly limited by their saturation temperature Tsat, which is as low as 100-150° K [7], [15], [16]. The other disadvantage of the SQUID-amplifiers is a narrow range of linearity of the transfer function. Implementation of a flux-locked-loop operating mode can substantially increase dynamic range and linearity, but at the same time the external feed-back loop will limit the maximum operation frequency to a few tens of megahertz at best. Therefore an internal negative feedback has been suggested in order to increase dynamic range and to linearize the transfer function of such an amplifier [8], [9]. However this is very problematic given typical low values of the SQUID-amplifier gain [7], [15], [16] of 10-15 dB, since higher amplification gain is needed to effectively achieve the negative feedback. What is needed is a way to use arrays of SQUIDs to achieve both greater dynamic range and greater linearity, without requiring such negative feedback. Linearity is particularly important for processing signals at radio frequencies, where a nonlinear transfer function can give rise to undesired harmonics and intermodulation products.\nOne approach to overcoming the drawbacks of SQUID-amplifiers is associated with multi-element Josephson structures and arrays, including Superconducting Quantum Interference Filters (SQIF). See Schopohl, U.S. Pat. No. 6,690,162, “Device for High Resolution Measurement of Magnetic Fields”; Oppenlander, U.S. Pat. No. 7,369,093, “Superconducting Quantum Antenna.” A SQIF is an array (parallel, series or parallel-series) of DC SQUIDs with an unconventional array structure [10-12]. The SQIF voltage response is characterized by a single sharp peak. Contrary to the usual SQUID, which shows unique properties due to the strict periodicity of its structure, the unique properties of the SQIF result from just the opposite, an unconventional non-periodic array structure. SQIFs are therefore a new development of an intelligent network of SQUIDs. SQIFs certainly offer an approach to achieving increased dynamic range, but this approach does not offer a clear way to achieve linearization of these fundamentally nonlinear devices."}
{"text": "The present invention pertains to a method of manufacture of fishing rods, and more particularly to a method of applying decorative windings to a rod shaft. Windings of the type to which this invention pertains are for decorative purposes and are placed along any part of a rod shaft which is free of hardware and are normally positioned just forward of the handle. The present invention is specifically concerned with intricate windings which, as will hereinafter become apparent, have heretofore been reserved for very expensive custom built rods.\nIn the past, intricate designs in decorative windings were produced by first tightly winding thread around the shaft to form a base for the decoration and then winding individual threads around the rod shaft over the base spiraling one thread in one direction and then spiraling the next in the opposite direction so that the threads intersect at precise intervals with the points of intersection forming a straight line parallel to the axis of the rod shaft. The process is repeated until a diamond, or like effect is achieved after which the windings are fixed and protected by application of clear epoxy or the like.\nThis process, while yielding a very acceptable product, suffers from numerous disadvantages from a commercial standpoint. More specifically, in order to produce a large or \"full-size\" diamond design with respect to the rod shaft diameter, with each diamond touching the next adjacent diamond, the exact position and angle must be determined, measured and laid out on the surface of the rod shaft prior to actual winding requiring a great deal of care and concentration. Each thread is then wrapped individually or in small groups side by side. After each wrap is made, the thread or groups of thread must be secured at both ends before proceeding with the next thread or group. Each of the foregoing steps is extremely time consuming and intricate and can usually be done only by the most skillful rod wrappers. This significantly increases the cost of the rod, thus prohibiting application of decorative windings to mass produced rods.\nIn addition to the cost factor, since the thread must be wound around the shaft, only spirals, cross-hatch or diamond designs have heretofore been possible. Large multi-colored cross-hatch or diamond designs are virtually impossible to produce uniformly, and to produce multi-colored designs, a thread for each color must be used which further complicates design and layout. Moreover, the known method is restricted to colors available for threads that can be used for fishing rods. Metallic or translucent threads are unavailable.\nThus, the only process presently available suffers from numerous problems and disadvantages and is completely inapplicable to rods for the average customer."}
{"text": "In recent years, as for a variable displacement oil pump, a two-stage discharge pressure characteristic is often required for supplying different apparatus and parts, whose required discharge pressures differ from each other, for example, moving engine parts and a variable valve actuation device configured to control engine-valve operating characteristics, with oil discharged from an oil pump. According to such a two-stage discharge pressure characteristic, the pump discharge pressure can be maintained at a first discharge pressure in a first pump speed range and also maintained at a second discharge pressure in a second pump speed range. One such variable displacement oil pump has been disclosed in Japanese Patent Provisional Publication No. 2008-52450 (hereinafter referred to as “JP2008-524500”), corresponding to International Publication No. WO 2006/066405 (A1).\nTo satisfy such a two-stage discharge pressure characteristic, the variable displacement oil pump, as disclosed in JP2008-524500, has a cam ring, which is moveable or pivotable against the spring force of a return spring. The variable displacement oil pump is configured to achieve the two-stage discharge pressure characteristic by supplying the discharge pressure (the pressurized working fluid) to a selected one of two pressure-receiving chambers defined on the outer peripheral surface of the cam ring and by changing an eccentricity of a geometric center of the cylinder bore of the cam ring with respect to the axis of rotation of a rotor (exactly, a vane rotor)"}
{"text": "In the optical communication industry, requirements of optical components are getting more and more severe. The optical devices suppliers have to provide smaller while cheaper components to meet customer's requirements. By applying semiconductor manufacturing process and other machining processes, kinds of high-precision and high-quality micro mechanical elements can be made in a mass production scale.\nMovable micro mirror chip is a kind of micro component applicable in optical communication or display modules. The micro mirror chip mainly includes a mirror layer and an actuation layer. The mirror layer is a suspended mirror actuated to swing or revolve by electrostatic force of the actuation layer. The micro mirror chip works as an attenuator in optical communication by steering the light beam direction. The biggest challenge in fabrication of a conventional movable micro mirror chip is the positioning alignment of the mirror layer and the actuation layer. The alignment process requires operations of hands or specific positioning machines that cause a high manufacturing cost very difficult to be reduced.\nU.S. Pat. No. 6,442,307 discloses a micro mirror device, in which the mirror layer and the actuation layer are jointed together by solder joints. Though the construction solves conduction and jointing problem, the solder joint does not provide precise positioning. Therefore, a spacer is user for the positioning function. However, when the solder melts, it loses its form. In order to stabilize the jointing and to remove the spacer after soldering, special materials for the spacer is required that makes the material selection very difficult. Also, the height of the solder layers is hard to be controlled after the large temperature variation. Moreover, since the mirror layer and the actuation layer all require electrical circuit to operate, there are wire-bonding areas on the layers that are easy to be stained by re-flow solder because there is no protection manner. The melted solder may stain the bonding areas and make the afterward wire bonding impossible, or even ruin the whole chip and increase fabrication cost of the product."}
{"text": "In recent years, lithium-ion secondary batteries and lithium-ion capacitors have been drawing attention as electric storage devices for use in mobile electronic devices such as mobile phones or laptop personal computers, electric vehicles, and hybrid vehicles. As anode current collectors for such electric storage devices, porous metal foils are used or are being considered for use. This is because making the foil porous provides benefits such that the weight of the foil can be reduced (to improve fuel consumption in automobiles), that adhesive power of an active material can be improved by anchoring effect making use of the pores, and that pre-doping of lithium ions (e.g., vertical pre-doping) can be efficiently conducted by making use of the pores.\nKnown methods for producing such porous metal foils include (1) a method of masking the surface of a substrate in a desired pattern with an insulating film, onto which electrolytic plating is conducted to form pores in accordance with the pattern; (2) a method of providing the surface of a substrate with a specific surface roughness or a specific surface condition, onto which electrolytic plating is conducted to control nucleation; (3) a method of perforating a non-porous metal foil by etching or machining; and (4) a method of forming a three-dimensional network structure by techniques for producing metal foams or plating nonwoven fabrics.\nIn particular, various techniques have been proposed for the above method of (2) since its steps are relatively simple and suitable for mass production. For example, Patent Literature 1 discloses a method for producing a fine-porous metal foil by subjecting an anode having a surface roughness Rz of 0.8 μm or less to electrolytic plating. Patent Literature 2 discloses a method comprising forming an oxidized film on the surface of a cathode body made of titanium or a titanium alloy by anode oxidization method; electro-depositing copper on the surface of the cathode body to form a porous copper foil; and peeling the foil from the cathode body. Patent Literature 3 discloses a method for producing a porous metal foil provided with an aluminum alloy carrier, comprising forming even projections by etching aluminum; and gradually growing metal particles from the projections as cores for electro-deposition so as to connect the metal particles.\nHowever, the actual situation is that the porous metal foils produced in accordance with these techniques fail to achieve a sufficiently high aperture ratio due to a concern that the foil strength would be lowered. In addition, a long foil is difficult to produce, and anode oxidation process had problems with the peelability of the porous metal foil and the stability of the aperture ratio, in that continuous peeling of the foil destroys the oxidized film. In particular, a higher aperture ratio is required for anode current collectors of electric storage devices, such as lithium-ion secondary batteries and lithium-ion capacitors, with improvement in performance."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display device and, in particular, to a reflection/transmission type liquid crystal display device capable of performing a display both in a reflection mode and a transmission mode.\n2. Description of the Related Art\nConventionally, there have been a reflection type liquid crystal display device utilizing ambient light, a transmission type liquid crystal display device utilizing backlight, and a semi-transmission type liquid crystal display device equipped with a half mirror and a backlight.\nIn a reflection type liquid crystal display device, a display becomes less visible under dim environment, whereas in a transmission type liquid crystal display device, a display becomes hazy under strong ambient light (e.g., under outdoor sunlight). As a liquid crystal display device capable of functioning in both modes so as to perform a satisfactory display under any environment, a semi-transmission type liquid crystal display device is disclosed by Japanese Laid-Open Publication No. 7-333598.\nHowever, the above-mentioned conventional semi-transmission type liquid crystal display device has the following problems.\nThe conventional semi-transmission type liquid crystal display device uses a half mirror in place of a reflective plate used in a reflection type liquid crystal display device, and has a minute transmission region (e.g., minute holes in a metal thin film) in a reflection region, thereby performing a display by utilizing transmitted light as well as reflected light. Since reflected light and transmitted light used for a display pass through the same liquid crystal layer, an optical path of reflected light becomes twice that of transmitted light, which causes a large difference in retardation of the liquid crystal layer with respect to reflected light and transmitted light. Thus, a satisfactory display cannot be obtained. Furthermore, a display in a reflection mode and a display in a transmission mode are superimposed on each other, so that the respective displays cannot be separately optimized. This results in difficulty in performing a color display, and causes a blurred display.\nA liquid crystal display device according to the present invention, includes: a first substrate, a second substrate, a liquid crystal layer interposed between the first substrate and the second substrate, and a plurality of pixel regions defined by a pair of electrodes for applying a voltage to the liquid crystal layer, wherein each of the plurality of pixel regions includes a reflection region and a transmission region, and the liquid crystal layer is made of a liquid crystal material having positive dielectric anisotropy, the device further including: a first polarizing element provided on the first substrate opposite to the liquid crystal layer; a second polarizing element provided on the second substrate opposite to the liquid crystal layer; a first phase difference compensator provided between the first polarizing element and the liquid crystal layer; and a second phase difference compensator provided between the second polarizing element and the liquid crystal layer, a twist angle xcex8t of the liquid crystal layer being in a range of 0xc2x0 to 90xc2x0, wherein retardation Rd and the twist angle xcex8t in a visible light region of the liquid crystal layer in the reflection region are in ranges within curves respectively represented by the following Formulae (1) and (2), and Formulae (3) and (4), in ranges within curves respectively represented by the following Formulae (5) and (6) and Formulae (7) and (8) at the twist angle xcex8t in a range of 0xc2x0xe2x89xa6xcex8txe2x89xa654.3xc2x0, and in ranges within curves respectively represented by the following Formulae (5) and (8) at the twist angle xcex8t in a range of 54.3xc2x0 less than xcex8txe2x89xa690xc2x0, and wherein the retardation Rd and the twist angle xcex8t in a visible light region of the liquid crystal layer in the transmission region are in ranges within curves respectively represented by the following Formulae (9) and (10) and Formulae (11) and (12):\nRd=xe2x88x920.0043xc2x7xcex8t2xe2x88x920.065xc2x7xcex8t+1011.8xe2x80x83xe2x80x83(1)\nRd=xe2x88x920.0089xc2x7xcex8t2+0.1379xc2x7xcex8t+914.68xe2x80x83xe2x80x83(2)\nRd=xe2x88x920.0015xc2x7xcex8t2xe2x88x920.1612xc2x7xcex8t+737.29xe2x80x83xe2x80x83(3)\nRd=xe2x88x920.0064xc2x7xcex8t2xe2x88x920.0043xc2x7xcex8t+640.65xe2x80x83xe2x80x83(4)\nRd=xe2x88x920.0178xc2x7xcex8t2+0.2219xc2x7xcex8t+458.92xe2x80x83xe2x80x83(5)\nRd=xe2x88x920.0405xc2x7xcex8t2+0.4045xc2x7xcex8t+364.05xe2x80x83xe2x80x83(6)\nRd=0.0347xc2x7xcex8t2xe2x88x920.4161xc2x7xcex8t+186.53xe2x80x83xe2x80x83(7)\nRd=0.0098xc2x7xcex8t2xe2x88x920.1912xc2x7xcex8t+89.873xe2x80x83xe2x80x83(8)\nxe2x80x83Rd=xe2x88x920.0043xc2x7xcex8t2xe2x88x920.065xc2x7xcex8t+995.66xe2x80x83xe2x80x83(9)\nRd=xe2x88x920.0058xc2x7xcex8t2xe2x88x920.0202xc2x7xcex8t+665.8xe2x80x83xe2x80x83(10)\nRd=xe2x88x920.0248xc2x7xcex8t2+0.6307xc2x7xcex8t+439.58xe2x80x83xe2x80x83(11)\nRd=0.0181xc2x7xcex8t2xe2x88x920.6662xc2x7xcex8t+109.51xe2x80x83xe2x80x83(12)\nIn one embodiment of the present invention, the retardation Rd is in a range within the curves respectively represented by Formulae (7) and (8) at the twist angle xcex8t in the reflection region in a range of 0xc2x0xe2x89xa6xcex8txe2x89xa654.3xc2x0, and in a range within the curves respectively represented by Formulae (5) and (8) at the twist angle xcex8t in the reflection region in a range of 54.3xc2x0 less than xcex8txe2x89xa690xc2x0, and the retardation is in a range within the curves respectively represented by Formulae (11) and (12) at the twist angle xcex8t in the transmission region in a range of 0xc2x0xe2x89xa6xcex8txe2x89xa690xc2x0.\nIn another embodiment of the present invention, the reflection region and the transmission region include a liquid crystal layer made of the same liquid crystal material, and a thickness of the liquid crystal layer in the reflection region is smaller than a thickness of the liquid crystal layer in the transmission region.\nIn another embodiment of the present invention, the first phase difference compensator has a first phase difference plate, the twist angle xcex8t of the liquid crystal layer is 0xc2x0, the retardation Rd of the reflection region is 90 nmxe2x89xa6Rdxe2x89xa6187 nm, the retardation Rd of the transmission region is 110 nmxe2x89xa6Rdxe2x89xa6440 nm, and the retardation Rd of the first phase difference plate is 30 nmxe2x89xa6Rdxe2x89xa6250 nm.\nIn another embodiment of the present invention, the first phase difference compensator further has a second phase difference plate, and the retardation Rd of the second phase difference plate is in a range of 220 nmxe2x89xa6Rdxe2x89xa6330 nm.\nIn another embodiment of the present invention, the second phase difference compensator has a third phase difference plate, and the retardation Rd of the third phase difference plate is in a range of 120xe2x89xa6Rdxe2x89xa6150 nm.\nIn another embodiment of the present invention, the second phase difference compensator further has a fourth phase difference plate, and the retardation Rd of the fourth phase difference plate is in a range of 240xe2x89xa6Rdxe2x89xa6310 nm.\nHereinafter, the function of the present invention will be described. First, the terms used herein will be described. In a reflection/transmission liquid crystal display device, a region where a display is performed by using transmitted light is referred to as a transmission region, and a region where a display is performed by using reflected light is referred to as a reflection region. The transmission region and the reflection region respectively include a transparent electrode region and a reflective electrode region formed on a substrate and a liquid crystal layer interposed between a pair of substrates. The transparent electrode region and the reflective electrode region on the substrate respectively define two-dimensional areas of the reflection region and the transmission region. The transparent electrode region is typically defined by a transparent electrode. The reflective electrode region is defined by a reflective electrode or a combination of the transparent electrode and the reflective electrode.\nThe liquid crystal display device of the present invention has a reflection region and a transmission region per pixel region. Thus, retardation of the liquid crystal layer can be optimized independently in the reflection region and the transmission region. More specifically, by prescribing the retardation of the liquid crystal layer in the reflection region to be those which (hatched regions (including double-hatched regions) in FIG. 5) are within curves represented by Formulae (1) and (2), Formulae (3) and (4), Formulae (5) and (6), and Formulae (7) and (8), and by prescribing the retardation of the liquid crystal layer in the transmission region to be those which (hatched regions (including double-hatched regions) in FIG. 6) are within curves represented by Formulae (9) and (10) and Formulae (11) and (12), the brightness (reflectivity) in the reflection region can be set to be about 70% or more, and the brightness (reflectivity) in the transmission region can be prescribed to be about 30% or more.\nIt is preferable that the conditions of the retardation are satisfied with respect to a central wavelength (high visibility) of visible light of about 550 nm. Furthermore, it is more preferable that the conditions of the retardation are satisfied in the entire wavelength range (about 400 nm to about 800 nm) of visible light.\nFurthermore, since the twist angle xcex8t is in a range of about 0xc2x0 to about 90xc2x0, the same twist angle can be obtained in both the reflection region and the transmission region having different thickenesses of the liquid crystal layer by single rubbing treatment. In order to render the twist angle different between the reflection region and the transmission region, rubbing is required to be conducted separately for two regions, which complicates a production process.\nFurthermore, by prescribing the retardation Rd in a region within the curves represented by Formulae (7) and (8) at the twist angle xcex8t of the reflection region in a range of 0xc2x0xe2x89xa6xcex8txe2x89xa654.3xc2x0, and in a region (double-hatched region in FIG. 5) within the curves represented by Formulae (5) and (8) at the twist angle xcex8t of the reflection region in a range of 54.3xc2x0xe2x89xa6xcex8txe2x89xa690xc2x0, and by prescribing the retardation Rd in a region (double-hatched region in FIG. 6) within the curves represented by Formulae (11) and (12) at the twist angle xcex8t of the transmission region in a range of 0xc2x0xe2x89xa6xcex8txe2x89xa690xc2x0, retardation of a liquid crystal layer in the reflection region and the transmission region becomes 0 in the presence of an applied voltage. If a black display is set to be performed at this time, a satisfactory black display is realized by applying the same voltage to the reflection region and the transmission region.\nFurthermore, the above-mentioned condition corresponds to the case where a white region in which retardation is closest to 0 (i.e., the first peak from the lowest retardation side in FIGS. 7 and 8) is selected as a condition of realizing a white display. Thus, a gray-scale display is also satisfactorily performed. More specifically, in a gray-scale state in which a white display is changed to a black display, brightness (reflectivity and transmissivity) is monotonously decreased, so that a satisfactory gray-scale display is obtained. If a white display is performed by using the second peak from the lowest retardation side in FIGS. 7 and 8, the first peak is present in a region for a gray-scale display. Thus, a satisfactory gray-scale display cannot be performed.\nWhen the liquid crystal layer in the transmission region and the reflection region are made of the same liquid crystal material, a structure and a production method will be simplified, compared with the case where the kind of a liquid crystal material is varied. It is effective to vary the thickness of the liquid crystal layer in the reflection region and the transmission region, so as to set different retardation in the reflection region and the transmission region. Furthermore, in order to match the length of an optical path with respect to light which contributes to a display in the reflection region with that in the transmission region, it is effective to prescribe the thickness of the liquid crystal layer in the transmission region to be larger than that in the reflection region. It is most preferable that the thickness of the liquid crystal layer in the transmission region is twice that in the reflection region.\nIf the first phase difference compensator has a first phase difference plate, the twist angle xcex8t of the liquid crystal layer is 0xc2x0, the retardation Rd of the reflection region is 90 nmxe2x89xa6Rdxe2x89xa6187 nm, the retardation Rd of the transmission region is 110 nmxe2x89xa6Rdxe2x89xa6440 nm, and the retardation Rd of the first phase difference plate is 30 nmxe2x89xa6Rdxe2x89xa6250 nm, a bright display of a normally white mode can be realized in the reflection region with a high contrast ratio.\nIf the first phase difference compensator has a second phase difference plate as well as the first phase difference plate, and the retardation Rd of the second phase difference plate is in a range of 220 nmxe2x89xa6Rdxe2x89xa6330 nm, wavelength characteristics in the reflection region can be alleviated, so that a display with a higher contrast can be obtained.\nIf the second phase difference compensator has a third phase difference plate, and the retardation Rd of the third phase difference plate is in a range of 120 nmxe2x89xa6Rdxe2x89xa6150 nm, a dark display is optimized even in the transmission region, so that a display with a higher contrast can be obtained.\nIf the second phase difference compensator has a fourth phase difference plate as well as the third phase difference plate, and the retardation Rd of the fourth phase difference plate is in a range of 240 nmxe2x89xa6Rdxe2x89xa6310 nm, wavelength characteristics of the transmission region are alleviated, so that a display with a higher contrast can be obtained.\nThus, the invention described herein makes possible the advantages of providing a liquid crystal display device which has outstanding mass-productivity and is capable of performing a satisfactory display irrespective of the brightness of ambient light.\nThese and other advantages of the present invention will become apparent to those skilled in the art upon reading and understanding the following detailed description with reference to the accompanying figures."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and related device for controlling operation of a portable electronic device, and more particularly, to a method and related device for determining whether a lid of a portable electronic device is open using a gravitational acceleration sensor and correspondingly determining the operation of the portable electronic device.\n2. Description of the Prior Art\nA laptop (i.e., a notebook computer) has several advantages, such as a small-sized volume, lightweight, and convenient for carrying due to its portability. These properties allow a user to work in any location. A small, thin, and light notebook computer provides the user with powerful computation abilities and document or multimedia processing functions anywhere and anytime, and thereby the work location of the user is not limited.\nPlease refer to FIG. 1. FIG. 1 is a schematic diagram of a notebook computer system 10 according to the prior art. Generally speaking, the notebook computer system 10 is composed of a lid 100 and a base 102. A hinge 104 connects the lid 100 and the base 102. The lid 100 comprises a screen, a camera, etc. The base 102 comprises a keyboard, a touchpad, a power switch, a host, an expanding interface, and so on. When using the notebook computer system 10, the user has to turn on the power of the host and adjusts a display angle of the screen of the lid 100 to a specific angle. In order to save power, a switch installed in the notebook computer system 10 can switch ON/OFF states of the screen and operation of the host according to an opening angle of the lid 100. For example, when the user doesn't need to use the notebook computer system 10 after turning on the notebook computer system 10, the user can close the lid 100 to make an angle, between the lid 100 and the base 102, smaller than a specific value, so that the notebook computer system 10 turns off the screen and operates in a sleep mode.\nAdjusting the angle between the lid 100 and the base 102, the user can save power and timely switch the operation of the notebook computer system 10. Therefore, it is considerably important to precisely detect the angle between the lid 100 and the base 102. In the prior art, there are many ways to detect the opening angle of the lid 100 and one of these is using a mechanic switch connected to the hinge 104. That is, turning off the screen and executing related operations, e.g. operating in the sleep mode when a rotating angle of the hinge 104 is smaller than a specific angle. However, the assembly of the mechanic switch is difficult and the mechanic switch may weary or malfunction by time, and finally, the reliability of the mechanic switch is decreased.\nIn addition, a magnetic sensor, such as a Hall sensor or a magnetic reluctance sensor, is used in the prior art. The notebook computer system 10 receives a distance from the lid 100 to the base 102 for determining the angle between the lid 100 and the base 102. For example, the Hall sensor can sense magnetic pole and magnetic force. Therefore, by installing a magnet in the lid 100 and a Hall sensor in the base 102, the notebook computer system 10 can determine the distance between the lid 100 and the base 102 so as to determine the angle between the lid 100 and the base 102. However, it is necessary to take the sensibility of the Hall sensor and magnetic flux of the magnet into account to meet demands when installing the magnet and the sensor. Besides, the magnetic reluctance sensor is difficult to design because of its high sensibility and narrow linear range."}
{"text": "The mining of gypsum or phosphogypsum produces incredible volumes of gypsum and phosphogypsum waste. This waste is typically stacked in very large piles, or gypsum stacks, which sometimes cover up to 400 acres in area to a depth of 100 to 200 feet or more. Thus, the land upon which the gypsum or phosphogypsum waste is stored is severely restricted with respect to any future uses. Therefore, gypsum stacks typically provide economic drain on the landowner or landholder.\nCurrently, a closed or capped phosphogypsum stack poses several economic challenges to the phosphate industry. In addition to the loss of use of the land upon which the stack is situate, governmental regulations that impose post-closure maintenance requirements necessitate ongoing expenses for decades. Until the current invention, gypsum stacks and other mineral waste sites, provided little or no economic return for gypsum mining companies to offset these ongoing expenses. Thus, for every company in the U.S. and abroad that mines, processes or produces gypsum, a gypsum stack always generates the inevitable: a large nonperforming asset. Numerous gypsum stacks currently exist with many more yet to come.\nThe present invention provides means and methods for converting presently existing gypsum stacks, or other mineral waste piles such as those found in any type of mining activity, into huge containers for holding refuse like biomass such as sugar, sugar cane, sugar cane waste, seaweed, fish, fish waste, shellfish, shellfish waste, agricultural waste, waste from forestry operations, other solid wastes, industrial waste or phosphate waste. The present invention is similarly applicable to other mineral waste piles such as those found in many types of mining activities, strip mining for example. The present invention also provides methods of doing business that will turn existing and future phosphogypsum stacks or mineral waste piles into income-producing waste containment and processing facilities.\nGypsum and waste phosphogypsum are capable of being formed into container structures which will withstand high compressive loads. Gypsum and phosphogypsum waste can therefore be used to form the bottom and walls of large containers. Thus, present gypsum stacks can be excavated to form large concavities which can be used as receptacles for solid waste. Current and future gypsum stacks can be formed into containers having one or more concavities as they are being created or enlarged. These concavities can be lined with geosynthetic liners or other membranes in order to prevent or limit leakage from the concavity. Examples of suitable geosynthetic fabrics, membranes, geocomposites, liners, liner combinations and liner-drain combinations and systems can be found in Designing With Geosynthetics, 2nd Edition, by Robert M. Koerner (1990, Prentice-Hall, ISBN 0-13-202300-8) which is incorporated herein by reference. Space in a concavity according to the present invention can be used to store waste materials, thereby providing an income-producing use of the gypsum stack or other mineral waste pile, and the land upon which they are located.\nThus, by providing hitherto unavailable waste storage facilities, the present gypsum stack and mineral waste pile conversion invention will significantly decrease the amount of new land that is used for landfills, while simultaneously converting gypsum stacks or waste piles from nonperforming assets into income-generating landfill facilities. The present means and methods for converting gypsum stacks or mineral waste piles into waste containment facilities is both economically and environmentally advantageous.\nFor instance, in comparison with conventional methods that would clear an additional 400 or more acres of forested or otherwise productive land to create a new MSW or C&D landfill, the present invention affords, among other advantages, at least one additional use for the presently existing tract of 400 or more acres of land by converting the space above, within or partially within an existing gypsum stack or mineral waste pile into waste containment or processing facilities. According to the present invention, instead of creating a new landfill or waste containment facility on a heretofore productive or uncleared parcel of land, the new landfill is constructed on top of, or wholly or partly within, a gypsum stack or other mineral waste pile, thereby sparing the productive land. Indeed, with the installation of one or more liner systems, and some minor modifications to the conversion site, the conversion site is very quickly able to receive waste such as biomass or other refuse.\nThe economic advantages of the present invention are significant. For instance, the governmental permits for establishing a 400-acre landfill are as much as $800,000.00. The present invention eliminates or substantially reduces such expenses because phosphogypsum stacks are already permitted as waste disposal sites. Thus, the present invention provides for an income-producing dual use of land upon which gypsum stacks or other mineral waste piles already exist.\nIn conventional use in the United States, when gypsum stacks reach maximum capacity at 100 to 200 feet of elevation, they are capped with an HDP membrane or liner in accordance with EPA regulations. In some embodiments of the present invention, the already existing gypsum stack cap membrane or liner can be used as a bottom liner for a waste containment concavity. Therefore, the cost to install a secondary liner in a landfill bottom liner system is reduced or eliminated in some embodiments of a phosphogypsum stack conversion according to the invention because the secondary liner is already in place in the form of the gypsum stack cap membrane or liner.\nAn additional economic advantage of the present invention pertains to subgrade costs. The costs relating to the preparation of the subgrade of a waste-containing site, which often requires the importation of clean subgrade material from offsite sources, can be as much as $2,000,000 U.S. for a 400-acre site. This expense is substantially reduced or eliminated with some embodiments of the present gypsum stack conversions since the phosphogypsum already available at the gypsum stack conversion site can be used as the subgrade material.\nOne significant aspect of the present invention is based upon the engineering parameters of mineral waste having high compressive strengths, such as waste gypsum, phosphogypsum or other mineral mining waste or byproducts. With the present invention, such waste or byproducts can be formed into one or more large concavities. Formation of a concavity according to the invention, can be effected, for instance, by excavation into an existing gypsum stack and then adding side walls of waste gypsum on top of the stack to extend the depth of the concavity to a desired dimension, by forming a concavity directly from phosphogypsum waste as the waste is being added to a particular geographic site, or by any other method or means that result in one or more concavities of desired dimensions. The concurrent delivery of waste and formation of the concavity is particularly useful when the gypsum waste is being transported to the site in the form of aqueous slurry, which may then be provided with strengthening binders or additives such as Portland cement as the concavity is being formed on top of an existing stack or as the concavity is being formed in place on the site. Thus, the cost of excavating a gypsum stack will be diminished.\nThe present invention can be used also with conventional means, devices and elements of the waste containment arts to thereby arrive at a facility of desired capabilities and capacities. For example, depending on the use to which a particular concavity is to be placed, liners or covers can also be provided to seal the bottom, walls or top of the concavity. Also in accordance with the invention, binders such as cement, for example Portland cement, can be combined with the phosphogypsum waste as it is being transported to the containment site. Addition of binders, such as cement, may be used to increase the compressive strength of the waste gypsum material and thereby assist in stabilizing the shape of the resulting concavity.\nThere is therefore a need for methods and means to provide additional and alternate uses for mineral waste piles, such as gypsum stacks. There is also a need for methods of forming phosphogypsum stacks into facilities useful for such purposes, and for business methods for producing revenue streams from such facilities."}
{"text": "The invention relates to a fuel supply unit comprising a pumping element and an electromotor arranged to drive the same. The pumping element and electromotor are housed in a common housing. A fuel supply unit is already known which is secured upright within a fuel container in such a manner that the pumping element is disposed near the bottom of the fuel container, a first intake area being provided oriented toward the bottom of the fuel container and a second intake area being provided which is remote from the bottom of the fuel container. As a result of this arrangement, it is possible to carry any vapor bubbles arising at elevated fuel temperatures away from the pump chamber and thus to assure satisfactory fuel supply. However, this arrangement also has the disadvantage that when the fuel level in the fuel container drops below the level of the second intake area, no further fuel is supplied. This stops the engine or prevents it from turning over upon starting, even though, since such containers at the present time are generally quite flat in shape, the container may still contain a relatively large quantity of fuel."}
{"text": "This invention relates to the field of separation science and analytical biochemistry.\nThe methods of this invention have applications in biology and medicine, including analysis of gene function, differential gene expression, protein discovery, cellular and clinical diagnostics and drug screening.\nCell function, both normal and pathologic, depends, in part, on the genes expressed by the cell (i.e., gene function). Gene expression has both qualitative and quantitative aspects. That is, cells may differ both in terms of the particular genes expressed and in terms of relative level of expression of the same gene. Differential gene expression can be manifested, for example, by differences in the expression of proteins encoded by the gene, or in post-translational modifications of expressed proteins. For example, proteins can be decorated with carbohydrates or phosphate groups, or they can be processed through peptide cleavage. Thus, at the biochemical level, a cell represents a complex mixture of organic biomolecules.\nOne goal of functional genomics (“proteomics”) is the identification and characterization of organic biomolecules that are differentially expressed between cell types. By comparing expression one can identify molecules that may be responsible for a particular pathologic activity of a cell. For example, identifying a protein that is expressed in cancer cells but not in normal cells is useful for diagnosis and, ultimately, for drug discovery and treatment of the pathology. Upon completion of the Human Genome Project, all the human genes will have been cloned, sequenced and organized in databases. In this “post-genome” world, the ability to identify differentially expressed proteins will lead, in turn, to the identification of the genes that encode them. Thus, the power of genetics can be brought to bear on problems of cell function.\nDifferential chemical analyses of gene expression and function require tools that can resolve the complex mixture of molecules in a cell, quantify them and identify them, even when present in trace amounts. However, the current tools of analytical chemistry for this purpose are limited in each of these areas. One popular biomolecular separation method is gel electrophoresis. Frequently, a first separation of proteins by isoelectric focusing in a gel is coupled with a second separation by sodium dodecyl sulfate-polyacrylamide gel electrophoresis (SDS-PAGE). The result is a map that resolves proteins according to the dimensions of isoelectric point (net charge) and size (i.e., mass). However useful, this method is limited in several ways. First, the method provides information only about two characteristics of a biomolecule—mass and isoelectric point (“pI”). Second, the resolution power in each of the dimensions is limited by the resolving power of the gel. For example, molecules whose mass differ by less than about 5% or less than about 0.5 pI are often difficult to resolve. Third, gels have limited loading capacity, and thus sensitivity; one may not be able to detect biomolecules that are expressed in small quantities. Fourth, small proteins and peptides with a molecular mass below about 10-20 kDa are not observed.\nOther analytical methods may overcome one or more of these limitations, but they are difficult to combine efficiently. For example, analytical chromatography can separate biomolecules based on a variety of analyte/adsorbent interactions, but multi-dimensional analysis is difficult and time consuming. Furthermore, the methods are limited in sensitivity.\nClinical diagnostics requires the ability to specifically detect known markers of disease. However, the development of such diagnostics is hampered by the time necessary to prepare reagents that specifically bind to markers, or that can discriminate the marker in a complex mixture.\nDrug discovery requires the ability to rapidly screen agents that modulate ligand/receptor interactions. Often the rate-limiting step in such screens is the ability to detect the ligand/receptor interaction. Thus, rapid and specific methods for identifying binding events would be an advance in the art.\nUntil now, the process from identifying a potential marker or member of a ligand/receptor pair to producing an agent that specifically binds the marker or member has been difficult. In one method, normal and diseased tissue are compared to identify mRNA species or expressed sequence tags (“ESTs”) that are elevated or decreased in the diseased tissue. These species are isolated and the polypeptides they encode are produced through routine methods of recombinant DNA. Then, the polypeptides are isolated and used as immunogens to raise antibodies specific for the marker. The antibodies can be used in, for example, ELISA assays to determine the amount of the marker in a patient sample.\nThis process is long and tedious. It can take nine months to a year to produce such antibodies, with much of the time being spent on developing protocols to isolate a sufficient quantity of the polypeptide for immunization. Furthermore, the method relies on the hope that differences in RNA expression are expressed as differences in protein expression. However, this assumption is not always reliable. Therefore, methods in which differentially expressed proteins are detected directly and in which specific ligands could be generated in significantly shorter time would be of great benefit to the field.\nThus, tools for resolving complex mixtures of organic biomolecules, identifying individual biomolecules in the mixture and identifying specific molecular recognition events involving one or more target analytes are desirable for analytical biochemistry, biology and medicine."}
{"text": "1. Reservoir Modeling\nSeismic data are routinely and effectively used to estimate the structure of reservoir bodies, but often play no role in the essential task of estimating the spatial distribution of reservoir properties. Reservoir property mapping is usually based solely on wellbore data, even when high resolution 3D seismic data are available. \"Wellbore data\" includes information typically obtained from a wireline log or core sample from an oil well and is typically the desired fine grain information for characterizing a \"reservoir property.\"\nPorosity, permeability, fluid and gas saturation, and other reservoir properties are measured at high accuracy near oil wells (e.g. by wireline logs), but these data do not assure reliable estimates of reservoir properties away from the wells. Seismic waves are not limited to wells, and seismic data may contain useful information about reservoir properties between the wells.\nProcessing of seismic data produces \"seismic attributes\" which may be effectively used to delineate the structure. See e.g., M. T. Taner, F. Koehler, and R. E. Sheriff, 1979, Complex Trace Analysis, Geophysics 44,1041-1063 and L. Sonneland, O. Barkved, and O. Hagness, 1990, Construction and Interpretation of Seismic Classifier Maps, EAEG meeting in Copenhagen. Seismic attributes are mathematical transformations on the data, computed either poststack (e.g., reflection intensity, instantaneous frequency, acoustic impedance, dip, azimuth) or prestack (e.g., AVO or moveout parameters). Well data are not used in their computation: attributes are purely properties of the seismic data alone. Otherwise, any analysis of the significance of seismic attributes to reservoir properties will be frustrated. Other seismic attributes include: acoustic impedance and velocity; reflection heterogeneity and instantaneous frequency; depth; dip and azimuth.\nSpecific seismic attributes may be related to specific reservoir properties. For example, acoustic impedance estimated from reflectivity by inversion of seismic data is an important seismic attribute. FIG. 13, shows cross sections of seismic data--reflectivity and acoustic impedance, together with wellbore data--porosity and water saturation. From FIG. 13, porosity does not appear directly related to the reflectivity, but it seems related to acoustic impedance--high impedance seems to imply low porosity. However, it is unclear how to actually use acoustic impedance to estimate porosity.\nOne problem in using seismic attributes is that their relation to rock properties is not obvious. For example, it is unclear how to use AVO to estimate gas saturation. Even if estimates are made, the confidence level in such estimates are unknown. There are unknown local factors that may affect the data in unexpected ways, and it is risky to predict functional relationships among seismic attributes and reservoir properties based on a simplified theoretical analysis with no familiarity of what \"works\" in a certain region. Region familiarity is built by comparing seismic and wellbore data. There is a need for interpretation methods and tools to build region familiarity, quantify its reliability, and subsequently use it to estimate properties. There is a need for a method to identify statistically-significant associations of seismic attributes and reservoir properties in any area, to determine the functional relationships implicit in these associations, to use them to predict the distribution of reservoir properties, and to quantify the reliability of the estimates.\n2. Artificial Neural Networks\nA great deal of recent research has been published relating to the application of artificial neural networks in a variety of contexts. See, for example, U.S. Pat. Nos. 5,134,685, 5,129,040, 5,113,483, and 5,107,442 (incorporated by reference). Artificial neural networks are computational models inspired by the architecture of the human brain. As a result three constraints are usually imposed on these models. The computations must be performed in parallel, the representation must be distributed, and the adjustment of network parameters (i.e., learning) must be adaptive. From an engineering perspective ANNs are adaptive, model-free estimators that estimate numerical functions using example data. While many different types of artificial neural networks exist, two common types are radial basis function (RBF) and back propagation artificial neural networks."}
{"text": "1. Field of the Invention\nThe present invention relates to an image pickup apparatus, such as a digital camera and a digital video camera.\n2. Description of the Related Art\nJapanese Patent Laid-Open No. (“JP”) 2003-295047 proposes a hybrid focus detection unit for moving a focus lens to a position near an on-focus position using an autofocus unit of a phase difference detection method (“phase difference AF”), and then for precisely moving the focus lens to the on-focus position using an autofocus unit of a contrast detection method (“contrast AF”). JP 2009-003122 realizes a phase difference detection function by providing an image pickup device with a focus detection pixel.\nHowever, the contrast AF has a problem in that a high contrast background in a focus detection region increases a contrast and causes an erroneous detection of this contrast change (which will be referred to as a “false peak” hereinafter) as an on-focus position. In addition, a noise ratio at the image pickup time increases as the ISO sensitivity increases, and the false peak is detected due to the noises. Hence, there is a demand for improved contrast AF accuracy or ultimately the hybrid AF accuracy."}
{"text": "Pulse oximetry systems for measuring constituents of circulating blood have gained rapid acceptance in a wide variety of medical applications, including surgical wards, intensive care and neonatal units, general wards, home care, physical training, and virtually all types of monitoring scenarios. A pulse oximetry system generally includes an optical sensor applied to a patient, a monitor for processing sensor signals and displaying results and a patient cable electrically interconnecting the sensor and the monitor. The monitor may be specific to pulse oximetry or may be a multi-parameter monitor that has a pulse oximetry plug-in. A pulse oximetry sensor has light emitting diodes (LEDs), typically one emitting a red wavelength and one emitting an infrared (IR) wavelength, and a photodiode detector. The emitters and detector are typically attached to a finger, and the patient cable transmits drive signals to these emitters from the monitor. The emitters respond to the drive signals to transmit light into the fleshy fingertip tissue. The detector generates a signal responsive to the emitted light after attenuation by pulsatile blood flow within the fingertip. The patient cable transmits the detector signal to the monitor, which processes the signal to provide a numerical readout of pulse oximetry parameters such as oxygen saturation (SpO2) and pulse rate."}
{"text": "This disclosure relates to an interfacial encapsulation process for making encapsulated beads particularly suited to enclose a fumigant effective against selected insect species. A typical insect species of concern is the fire ant. Fire ants are difficult to kill with ingested poisons. While ingested poisons may thoroughly decimate foraging worker ants, the use of a chain of tasters in the fire ant colony prevents ingested poison from reaching the queen of the colony, thereby protecting the colony. A particularly valuable volatile insecticidally effective poison includes volatile phosphoric or thiophosphoric acid esters. In addition to operating by ingestion, they provide a fumigant which is airborne, by-passing the tasters. Thus, if the bait is carried into the colony and opened in that closed environment, there is a much greater possibility that the bait will be effective to kill the queen of the colony. The fumigant is is commonly known as DDVP, a fumigant including as one ingredient dimethyl-2, 2-dichlorovinyl phosphate acid ester. There are other fumigant phosphoric or thiophosphoric acid esters which will be collectively referred to also as DDVP. As an example, several such insecticidally effective compounds or mixtures thereof are set forth in U.S. Pat. No. 4,094,970.\nDDVP produces a toxic atmosphere for the ants in the colony particularly if the bait can be carried into the colony by foraging ants and is opened in that closed atmosphere. The ants are enticed to open it by incorporation of an attractant food comprising the shell. One suitable attractant food is soybean protein. Hence, the present apparatus contemplates fabrication of a small bead by means of microencapsulation wherein the shell is fabricated with a small measure of soybean protein, and the core is an edible mixed with DDVP. The amount of DDVP is sufficient to be effective in the atmosphere; that is, once the bait has been broken open, DDVP volatilizes sufficiently to have the fatal impact required for colony decimation.\nThe bead shell is an important factor in preparing the bead of bait material for eventual extermination of a fire ant colony. On the one hand, it should be frangible and broken easily by the insect so that the insect can readily bite through the surrounding shell. On the other hand, the shell must be impervious to water from the exterior, impervious to the core on the interior to prevent DDVP from weeping or saturating the shell material whereby premature fumigation might occur. If this were to happen, premature toxicant exposure might well warn away the foraging ants. Moreover, toxicant exposure might well poison foraging ants before they have the opportunity to deliver the bead into the colony. Thus, it is desirable that the ant finding the bead remove it as a result of foraging; this is a delicate balance wherein the bead must be tough enough to be handled, delivered by machine, and yet should be sufficiently easily opened by insect bite. Moreover, it should be sufficiently free of DDVP which might possibly weep through the wall of the shell and thereby defeat the highly desirable fumigant procedure for insect eradication.\nOne prior art structure is found in U.S. Pat. No. 4,094,970 setting forth a polyurethane system. Additional references are U.S. Pat. Nos. 3,492,380; 3,575,882; 3,270,100 and 3,577,515. By and large, they generally refer to solid carrier pelletization processes. It is submitted, however, that a quality insecticidal bait must have a toxicant which is enclosed within the bead formed by interfacial encapsulation and there should be an attractant food dissolved in the shell. That is, the attractant food material must in some fashion be in the shell to attract foraging insects. Otherwise, the foraging insects will have no interest in the bead and will ignore it.\nThe core in the bead is made with DDVP as the fumigant mixed with soybean oil. Separately, the shell ultimately formed should have an attractant such as soybean protein or oil in or on the shell to serve as an attractant. It must either taste or smell good to the foraging insect, sufficiently to cause the insect to carry the bead back to the colony or hive. In this light, it will be understood that the completed bead has the secure, impervious surrounding shell which encapsulates the DDVP. Fumigation does not start until the shell is actually broken open. In the meanwhile, the shell is sufficiently attractive to the foraging insect that it will be carried back to the hive or colony.\nInterfacial polycondensation encapsulation is one procedure which enables this to be accomplished. As an example, an interfacial liquid body is defined as having two portions, one being an oil phase which forms droplets in water. The oil phase is added to the water to react with the oil phase at the droplet interface. The polycondensate system is defined by the two parts, one part being the oil phase and the other being the water phase. This method brings the two parts together, thereby achieving bead formation at the droplet interfacial area. With controlled stirring and the addition of either a surfactant or anti-foaming agent, bead size is controlled to form coated beads which are recovered and washed. By this procedure, beads of a selected size are formed and are recovered, washed and thereafter used. The interfacial polycondensate system obtains production of beads in the range of about 1.0 mm or smaller wherein the surrounding shell is impervious to DDVP enclosed therein."}
{"text": "Sausage making machines such as that shown in U.S. Pat. No. 3,115,668 include a meat emulsion pump connected to an elongated stuffing tube which extrudes into a shirred casing thereon a strand of sausage into the rotatable chuck of a twister housing. The strand then moves through a conventional linker and the linked strand discharged therefrom typically moves into a rotatable looper horn and deposits on the hooks of a moveable conveyor.\nIt would be advantageous to use such machines in spaced parallel relationship, to be attended by a single operator. However, conventional sausage making machines do not make this possible because the chuck and the looper horn can only be rotated in a single direction.\nIt is therefore the principal object of this invention to provide a sausage making machine which has power means associated with both the twister housing and the looper horn to permit the chuck in the twister housing and/or the looper horn to be rotated in both clockwise or counterclockwise directions."}
{"text": "In conventional speaker systems, there are solutions for controlling individual speakers or using a control component for managing a group of speakers. However, these conventional solutions rely upon wired connections or, in the case of wireless connections, individual speakers are often controlled by a single device, which is often inflexible and confines media to that selected using the single control device. Further, conventional solutions are often time-consuming and technically complex to set up and manage, often requiring extensive training or expertise to operate.\nConventional media playback solutions are typically found in mobile devices such as mobile phones, smart phones, or other devices. Unfortunately, conventional speaker control devices are often limited connections between a mobile device and a single speaker. Further, the range of actions that can be taken are often limited to the device that is in data communication with a given speaker. If different users with different playlists and mobile devices want to use a given speaker, individual connections often need to be established manually regardless of the type of data communication protocol used.\nCurrent radio standards (e.g., Bluetooth systems, WiFi systems) allow for a receiver to measure signal strength (e.g., of a RF signal) from a source transmitting data and one measure of signal strength includes received signal strength (RSSI). Although there have been studies that utilize RSSI information to understand how well RSSI values correlate to how far away a transmitter and a receiver are from one another, it is also known that it is difficult to utilize RSSI for distance measurements due to a number of factors. One of those factors may include a multipath effect where the RF signal being transmitted reflects off of surrounding objects, such as walls, stationary objects, and moving objects. Another factor may include antenna radiation pattern and polarization of antenna of the transmitter and the antenna of the receiver, both of which may contribute to RSSI error vs. distance. However, close distance measurements perform with higher accuracy than long distance measurement due to an inverse square power drop off (e.g., 1/R2 where R=Distance) in a far field region, and where for a near field region the inverse power drop can be greater than 1/R3 of the RF signal as a function of distance between the transmitter and the receiver. Close proximity sensing can be utilized to improve intuitiveness on how two or more devices interact with one another rather than having a user interact with them. One example is for the user to place one of the devices close to another device, within boundaries of a set threshold RSSI for close proximity detection. Although close proximity sensing via RSSI may have a statistically high level of accuracy and a device may infer that two devices are close to one another, there still exists a small probability that a false alarm can be triggered (i.e., the device is detected as being in close proximity, but actually in reality the device is not in close proximity). In conventional implementations, use cases would require perfect or near perfect inference of close proximity of the devices.\nThus, a need exists for a for speaker control solution without the limitations of conventional techniques and a solution that does not trigger false alarms when a received RSSI value is within a pre-determined RSSI threshold value, but the devices are not within close proximity of one another.\nAlthough the above-described drawings depict various examples of the present application, the present application is not limited by the depicted examples. It is to be understood that, in the drawings, like reference numerals designate like structural elements. Also, it is understood that the drawings are not necessarily to scale."}
{"text": "The present invention relates to a facsimile machine and, more specifically, relates to a facsimile machine that can provide a charge management system in which beneficiaries are charged costs of expendable supplies in accordance with the amounts they have spent.\nConventionally, where one facsimile machine is commonly used by a plurality of business sections, there has been a problem of how to share the related expenses. There are known the following patent documents which describe techniques for sharing such expenses.\n(1) Japanese Unexamined Patent Publication No. Sho. 56-10773: A key counter is activated when document information is received by a facsimile machine. Since the key counter has functions of starting the facsimile machine and accumulating the charge, the communication costs can be shared by beneficiaries.\nThere also exists a technique which employs a magnetic card as a key. Copying and transmission operations are prohibited unless the card is set into the machine. If the card is properly set, the number of copied sheets or the communication charge are registered in association with a business section number recorded in the card being set into the machine.\n(2) Japanese Unexamined Patent Publication No. Sho. 58-133073: A sender is charged for the communication time of image information and for the quantity of recording sheets he has spent. Where a direct mail is sent through a facsimile machine, a sender is charged for recording sheets consumed.\n(3) Japanese Unexamined Patent Publication No. Sho. 62-227262: Communication costs are distributed to users in a facsimile machine connected to an ISDN network.\nExcept for special cases of, e.g., the direct mail as described in above item (2), it would be appropriate that when information is received by a facsimile machine, costs of expendable supplies, such as recording sheets, of the receiving-side facsimile machine be paid by a receiver. However, in the prior art as described above, the costs for the expendable supplies cannot be distributed to receivers. That is, the expenditures of a facsimile machine cannot be shared in a correct manner."}
{"text": "Traditional mainframe computer configurations provided for user interface to the computer through computer terminals which were directly connected by wires to ports of the mainframe computer. As computing technology has evolved, processing power has typically evolved from a central processing center with a number of relatively low-processing power terminals to a distributed environment of networked processors. Examples of this shift in processing include local or wide area computer networks which interconnect individual work stations where each workstation has substantial independent processing capabilities. This shift may be further seen in the popularity of the Internet which interconnects many processors and networks of processors through devices such as, for example, routers. This type of network environment is often referred to as a client-server environment with client stations coupled to and supported by a server station.\nIn the modern distributed processing computer environment, control over software, such as application programs, is more difficult than where a mainframe operated by an administrator is used, particularly for large organizations with numerous client stations and servers distributed widely geographically and utilized by a large number of users. Furthermore, individual users may move from location to location and need to access the network from different client stations at different times. The networked environment increases the challenges for a network administrator in maintaining proper licenses for existing software and deploying new or updated application programs across the network.\nA further complication in network systems is that, typically, these systems include combinations of network applications and native applications as well as combinations of different connection types and hardware devices. As used herein “native applications” refers to applications which are installed locally on a workstation such that characteristics associated with the native application are stored on the workstation. The combinations of network connections, differing hardware, native applications and network applications makes portability of preferences or operating environment characteristics which provide consistency from workstation to workstation difficult. Furthermore, differences in hardware or connections may create inefficiencies as users move from workstation to workstation. For example, a user may, in a first session, access the network utilizing a high speed connection and a workstation with a high resolution color monitor to execute an application and then, in a later session, access the network to execute the same application from a mobile computer with a monochrome display and a low speed modem connection to the network. Thus, session content, such as color display data or preferences associated with the application, which may have been appropriate for the first session may be inappropriate or inefficient in a later session.\nEfforts to address mobility of uses in a network have included efforts to provide preference mobility such as, for example, Novell's Z.E.N.works™, Microsoft's “Zero Administration” initiative for Windows® and International Business Machines Corporation's (IBM's) Workspace On Demand™. However, these solutions each typically require pre-installation of software at the workstation to support their services. For example, Novell's Z.E.N. and IBM's Workspace On Demand utilize a vendor-supplied support layer in the operating system to enable their services. In addition to modifying the workstations operating system at startup to setup tasks to customize the user's environment, the Microsoft Zero Administration solution may be limited to a homogeneous environment where the workstation and the server are utilizing the same operating system.\nEach of these “mobility” systems typically do not address the full range of complications which may arise in a heterogeneous network utilizing differing devices and connections. Users would typically have to manually define session characteristics at each differing workstation they used in the network or maintain local characteristic definitions which may be inappropriate for particular applications a user is executing and may substantially reduce the administrative convenience of a centrally controlled network. Thus, these various approaches fail to provide a seamless integration of session characteristics across heterogeneous network devices. Such solutions may reduce network administration only after initial installation on each workstation. Furthermore, content is typically not addressed such that inefficiencies in use of the network may result."}
{"text": "This invention relates generally to computer architectures and more particularly to memory and busing architectures within computers.\nComputers are known to include a central processing unit, system memory, a memory controller, a chip set, video graphics circuitry, interconnecting buses, and peripheral ports. The peripheral ports enable the central processing unit and/or other components, to communicate with peripheral devices such as monitors, printers, external memory, etc.\nIn most computer systems, a computer will include cache memory to more effectively access larger memory, such as a system memory hard drive. As is known, cache memory is relatively small in comparison to system memory and can be accessed by the central processing unit much more quickly than the system memory. As such, when the central processing unit has a read and/or write function to process pertaining to a particular data element stored in the system memory, the data element, and related data elements, are retrieved from system memory and provided to the cache memory. As is also known, the related data elements may be in the same memory line (e.g., 128 bytes) as the data element, or in the same memory block (e.g., several memory lines) as the data element.\nThe rationale for retrieving a line or several lines of memory is based on the assumption that the central processing unit is processing sequentially related operations that have data elements stored in groupings (i.e., in the same memory line or group of memory lines). For example, video graphics data is often stored in a linear or tiled manner, wherein the memory locations correspond to pixel locations of the display. As such, filling the cache with the needed data element and the related data elements requires only one read operation from the system memory, while the central processing unit may perform multiple reads and/or writes of data in the cache. Thus, memory access is much more efficient. Note that the retrieved data elements may be related temporally or spatially.\nAs is also known, data elements are bused in a pipeline manner wherein, for a given transaction, i.e., a read and/or write of data, the transaction includes an arbitration phase, a request phase, an error phase, a response phase, and a data phase. Each of these phases may be several clock cycles in length and their length varies depending on the busyness of the bus. As such, some data elements may be processed quickly while others are delayed or preempted due to higher priority data traffic. As is known, the error phase is used to determine whether a particular transaction is being preempted for a higher priority transaction. Thus, additional processing is required due to the varying processing lengths of transactions.\nIn some computer systems, there are multiple processors, where each processor has its own cache, which may include two levels of cache. The first level cache being smaller and more readily accessible than the second level. Thus, in such multiprocessor environments, when a processor is not utilizing its cache, it remains idle. Conversely, when a processor is processing a significant amount of data its cache may be too small, thus forcing the data to be thrashed between cache memory and system memory, which is inefficient. As such, the cache memory in a multiprocessor environment is not used as effectively as possible.\nTherefore, a need exists for a method and apparatus that more efficiently utilizes cache memory and more efficiently buses data within a computer system."}
{"text": "The use of Radio Frequency Identification (RFID) tags for inventory control is well known. These tags are broadly defined as radio frequency transponders allow tagged inventory items to uniquely self identify themselves to a suitably configured network of RFID readers. These readers interface with a computer network to monitor the movement and/or status of inventory or work-in-process.\nSystems tend to fall into two categories, fixed point monitoring and scanning. In the fixed point modality tagged items move passed a fixed reader that localizes the item in time and position. This modality is well suited to monitoring movement of items into and out of a specific area as well as movement along a linear process such as a packaging or assembly line. In a scanning modality, a mobile scanner moves in an area reading and logging, all the RFID tags that are in range of the reader as it moves through the inventory area. This modality is better suited for monitoring static inventory such as materials in a warehouse or books in an archive, in further discussion, it should be understood that warehouse may be taken to mean for holding static inventory.\nIn many warehouse applications using a scanning or roaming reader, it is also advantageous to know not only that an item is present in the warehouse but also where the item is located, i.e. on which shelf our in which aisle in a storeroom. In order to localize items in a space, the roaming reader may follow a preprogrammed track or may periodically update its location through a number of methods. The mobile reader may have an integrated Real Time Location System (RTLS), the reader may have its location manually updated by an operator, or it may pass certain reference points in its progress along a preplanned route through the warehouse. In general these approaches are adequate for determining the inventory state of a warehouse; however they have the disadvantage that the reader must follow a pre-programmed (deterministic) path or route through the warehouse and they require an initial sweep to determine the initial locations of items in the warehouse. This may be time consuming and may not provide the degree of real time location of items that is required, especially in an active warehouse where items may be moved from storage location to storage location or there may be several stocking locations for similar or associated items with the same warehouse. Further the ease of monitoring, storing and moving data associated with a large operation may be difficult."}
{"text": "The invention pertains to the practice of cleaning objects. The invention solves the problems of lack of abrasion and lack of durability found in prior swabs (such as cotton tipped or foam tipped)."}
{"text": "The present invention relates to corona treaters and more specifically to the electrical connection between the electrode of the corona treater and a high voltage source. A corona treater utilizes an electrode magazine or cassette that must be periodically removed from the corona treater. In the past, the electrical connection between the high voltage source and the electrode magazine was a permanently soldered connection that made removal of the electrode magazine from its station a time consuming and somewhat difficult job.\nThe purpose of the present invention is to provide a simple high voltage plug-in connection between the electrode and the high voltage source so that the connection can be easily disconnected to facilitate the removal of the electrode magazine from its station."}
{"text": "This invention relates to preparing fluorinated electrets.\nThe filtration properties of nonwoven polymeric fibrous webs can be improved by transforming the web into an electret, i.e., a dielectric material exhibiting a quasi-permanent electrical charge. Electrets are effective in enhancing particle capture in aerosol filters. Electrets are useful in a variety of devices including, e.g., air filters, face masks, and respirators, and as electrostatic elements in electro-acoustic devices such as microphones, headphones, and electrostatic recorders.\nElectrets are currently produced by a variety of methods including direct current (xe2x80x9cDCxe2x80x9d) corona charging (see, e.g., U.S. Pat. Re. 30,782 (van Turnhout)), and hydrocharging (see, e.g., U.S. Pat. No. 5,496,507 (Angadjivand et al.)), and can be improved by incorporating fluorochemicals into the melt used to produce the fibers of some electrets (see, e.g., U.S. Pat. No. 5,025,052 (Crater et al.)).\nMany of the particles and contaminants with which electret filters come into contact interfere with the filtering capabilities of the webs. Liquid aerosols, for example, particularly oily aerosols, tend to cause electret filters to lose their electret enhanced filtering efficiency (see, e.g., U.S. Pat. No. 5,411,576 (Jones et al.)).\nNumerous methods have been developed to compensate for loss of filtering efficiency. One method includes increasing the amount of the nonwoven polymeric web in the electret filter by adding layers of web or increasing the thickness of the electret filter. The additional web, however, increases the breathing resistance of the electret filter, adds weight and bulk to the electret filter, and increases the cost of the electret filter. Another method for improving an electret filter\"\"s resistance to oily aerosols includes forming the electret filter from resins that include melt processable fluorochemical additives such as fluorochemical oxazolidinones, fluorochemical piperazines, and perfluorinated alkanes. (See, e.g., U.S. Pat. No. 5,025,052 (Crater et al.)). The fluorochemicals should be melt processable, i.e., suffer substantially no degradation under the melt processing conditions used to form the microfibers that are used in the fibrous webs of some electrets. (See, e.g., WO 97/07272 (Minnesota Mining and Manufacturing)).\nIn one aspect, the invention features an electret that includes a surface modified polymeric article having surface fluorination produced by fluorinating a polymeric article. In one embodiment, the article includes at least about 45 atomic % fluorine as detected by ESCA. In another embodiment, the article includes a CF3:CF2 ratio of at least about 0.25 as determined according to the Method for Determining CF3:CF2. In other embodiments, the article includes a CF3:CF2 ratio of at least about 0.45 as determined according to the Method for Determining CF3:CF2.\nIn one embodiment, the article has a Quality Factor of at least about 0.25/mmH2O, (preferably at least about 0.5/mmH2O, more preferably at least about 1/mmH2O).\nIn some embodiments, the article includes a nonwoven polymeric fibrous web. Examples of suitable fibers for the nonwoven polymeric fibrous web include polycarbonate, polyolefin, polyester, halogenated polyvinyl, polystyrene, and combinations thereof. Particularly useful fibers include polypropylene, poly-(4-methyl-1-pentene), and combinations thereof. In one embodiment, the article includes meltblown microfibers.\nIn another aspect, the invention features an electret that includes a polymeric article having at least about 45 atomic % fluorine as detected by ESCA, and a CF3:CF2 ratio of at least about 0.45 as determined according to the Method for Determining CF3:CF2. In another embodiment, the electret includes at least about 50 atomic % fluorine as detected by ESCA, and a CF3:CF2 ratio of at least about 0.25 as determined according to the Method for Determining CF3:CF2.\nIn other aspects, the invention features a respirator that includes the above-described electrets. In still other aspects, the invention features a filter that includes the above-described electrets.\nIn one aspect, the invention features a method of making an electret that includes: (a) fluorinating a polymeric article to produce an article having surface fluorination; and (b) charging the fluorinated article in a manner sufficient to produce an electret. In one embodiment, the method includes charging the fluorinated article by contacting the fluorinated article with water in a manner sufficient to produce an electret, and drying the article. The method is useful for making the above-described electrets. In another embodiment, the method includes charging the fluorinated article by impinging jets of water or a stream of water droplets onto the fluorinated article at a pressure and for a period sufficient to produce an electret, and drying the article.\nIn other embodiments, the method includes fluorinating a polymeric article in the presence of an electrical discharge (e.g., an alternating current corona discharge at atmospheric pressure) to produce a fluorinated article. In one embodiment, the method includes fluorinating the polymeric article in an atmosphere that includes fluorine containing species selected from the group consisting of elemental fluorine, fluorocarbons, hydrofluorocarbons, fluorinated sulfur, fluorinated nitrogen and combinations thereof. Examples of suitable fluorine containing species include C5F12, C2F6, CF4, hexafluoropropylene, SF6, NF3, and combinations thereof.\nIn other embodiments, the method includes fluorinating the polymeric article in an atmosphere that includes elemental fluorine.\nIn other embodiments, the method of making the electret includes: (A) fluorinating a nonwoven polymeric fibrous web (i) in an atmosphere that includes fluorine containing species and an inert gas, and (ii) in the presence of an electrical discharge to produce a web having surface fluorination; and (B) charging the fluorinated web in a manner sufficient to produce an electret.\nIn other aspects, the invention features a method of filtering that includes passing an aerosol through the above-described electrets to remove contaminants.\nThe fluorinated electrets of the invention exhibit a relatively high oily mist resistance relative to non-fluorinated electrets.\nIn reference to the invention, these terms having the meanings set forth below:\nxe2x80x9celectretxe2x80x9d means a dielectric material exhibiting a quasi-permanent electrical charge. The term xe2x80x9cquasi-permanentxe2x80x9d means that the time constants characteristic for the decay of the charge are much longer than the time period over which the electret is used;\nxe2x80x9csurface modifiedxe2x80x9d means that the chemical structure at the surface has been altered from its original state.\nxe2x80x9csurface fluorinationxe2x80x9d means the presence of fluorine atoms on a surface (e.g., the surface of an article);\nxe2x80x9cfluorine containing speciesxe2x80x9d means molecules and moieties containing fluorine atoms including, e.g., fluorine atoms, elemental fluorine, and fluorine containing radicals;\nxe2x80x9cfluorinatingxe2x80x9d means placing fluorine atoms on the surface of an article by transferring fluorine containing species from a gaseous phase to the article by chemical reaction, sorption, condensation, or other suitable means;\nxe2x80x9caerosolxe2x80x9d means a gas that contains suspended particles in solid or liquid form; and\nxe2x80x9ccontaminantsxe2x80x9d means particles and/or other substances that generally may not be considered to be particles (e.g., organic vapors)."}
{"text": "This application relates to video input circuits and, more particularly, to video input circuits for a video hard copy controller of the type that converts and formats video signals into digital signals for application to a hard copy generating device.\nPrior art video input circuits have included a video amplifier circuit having input terminals for receiving a composite video signal comprised of a video data component and a synchronization component, control terminals for receiving a gain control signal to control the gain of the amplifier circuit, and output terminals at which an amplified composite video signal is developed.\nIn the past, a combination automatic gain control and d-c restorer circuit was also employed to generate the gain control signal and to restore the tip extremity of the synchronization component to a reference potential, such as a d-c common.\nA major problem with the above video input circuits arises from the fact that the amplified composite video signal at the output of the video amplifier was first started to be d-c restored and then, during the d-c restoration process, the automatic gain control (AGC) circuit was enabled. Although d-c restoration started before the AGC circuit was enabled and the AGC circuit completed its first operational iteration after d-c restoration was terminated, yet there was a significant period of time during which the d-c restoration and AGC operations overlapped. This posed significant conflicting effects with unreliable results. Also, since the tips of the synchronization component were restored to the d-c reference common, there was substantial likelihood of unwanted ground noise influence.\nIt would be desirable to provide video input circuits of the general type above described wherein the undesirable overlap of d-c restoration and AGC operations is eliminated. It would further be desirable if such video input circuits could include a simpler and more efficient d-c restorer circuit that could reduce ground noise."}
{"text": "The invention relates to an electro-surgical treatment instrument comprising a holder connectable to a high frequency generator with at least two electrodes provided at the proximal end of the holder which can be brought into contact with the tissue of a patient.\nIn a known treatment instrument of this kind, (DE 34 23 352 C2) the two electrodes are arranged alongside one another with respect to the axis of the treatment instrument and formed as the branches of a pair of tweezers. In this way, tissue regions of a patient which can be grasped in a tweezer-like manner with the electrodes can be subjected to bi-polar coagulation.\nIf, however, tissue regions are to be coagulated which cannot simply be grasped in a simple tweezer-like manner, then coagulation requires very much attention and skill on behalf of the operator. This is because he must then select the electrode spacing necessary for bi-polar coagulation by pressing the tweezer-like electrodes together and adapt it as necessary to the resulting coagulation."}
{"text": "In fabricating semiconductor devices, packaging is a process to protect the semi-conductor chips from an external environment, to shape a semiconductor chip for an easy application, and to protect the operation functions added in the semiconductor chip and thus enhance the reliability of the semiconductor device.\nRecently, as the semiconductor devices are highly integrated and their functions become versatile, the packaging is gradually transferred from a process having a small number of package pins to a process having a large number of package pins, and is also converted from a structure where a package is inserted into a printed circuit board (PCB) to a structure where a package is mounted on a surface of a PCB, i.e., surface mounting device structure. Examples of the surface mounting type packages include a small outline package (SOP), a plastic leaded chip carrier (PLCC), a quad flat package (QFP), a ball grid array (BGA), a chip scale package (CSP) and the like.\nA chip carrier related with these semiconductor packages or a base substrate used in a PCB should be stable thermally, electrically and mechanically. As the chip carrier or the base substrate for a PCB, a high price ceramic substrate or a resin substrate having a polyimide resin, a fluorine resin or a silicon resin as a base material has been used.\nSince the ceramic substrate or resin substrate is made of insulator, it need not deposit an insulation material after the through hole process. However, in the case of resin substrates, since their material cost is expensive and is poor in the moisture-resistant property and heat-resistant property, it is not good to use the resin substrates as the chip carrier substrate. Also, although the ceramic substrate is comparatively superior in terms of heat-resistant property to the resin substrate, the ceramic substrate is also expensive and has disadvantages of a high processing cost as well as a difficulty in the processing.\nTo overcome the disadvantages of these ceramic substrates or resin substrates, use of a metallic substrate is proposed. The metallic substrate is advantageous in that it is inexpensive, can be easily processed and has a good thermal reliability. However, the metallic substrate requires a separate insulation treatment, which is unnecessary for the aforementioned ceramic or resin substrates, and also requires to attach a metal core serving as a heat sink or a heat spread on an upper surface or a lower surface of a completed substrate so as to more effectively irradiate heat.\nIn the meanwhile, the chip carrier or PCB prefers a thin and flat one in accordance with the current design trend toward a lightweight, slim and miniaturized profile. To realize the slimness and flatness, a technique that a cavity is formed in a substrate and a chip or a component is mounted on the formed cavity is employed.\nIn the case of resin substrates, such a cavity is formed by drilling the resin substrate. However, the drilling method takes much time and high cost in processing the cavity. Also, the cavities as formed may have a large deviation, which allows a mounted component to be leaning to one side and makes it difficult to maintain the flatness. Furthermore, since the resin used as the material of the substrate is poor in thermal and mechanical characteristics, when a component is mounted on the substrate, a serious deformation may be caused due to a stress."}
{"text": "Cable construction and central office re-concentration or replacement projects often require half-tapped or double-tapped placement of wire or cabling to facilitate the conversion or construction process. As part of the process of installing new cables and/or equipment, bridge taps often occur. A bridge tap is a length of wire or cable attached to normal endpoints of a circuit that introduces unwanted impedance imbalances that can interfere with data transmission.\nIn some instances, cables may be placed weeks or months in advance of the actual conversion or cable use. When pre-run, the cable ends are often laid in place, with one end of the new cable being connected to an existing cable and the other end left unterminated resulting in a bridge tap.\nBridge tap causes a wire to reflect signals from the unterminated end back to the source. Data signals, e.g., DSL signals, operate at high frequencies and can be severely impacted by the presence of bridge tap reflections. Circuits may operate at lowered speeds or data rates as a result of the bridge tap interference. Marginal circuits could exceed the operational limits of a design, thus entirely preventing the circuit from operating.\nAccordingly, bridge taps can cause problems in operating systems. As the quality of service and attainable data rates degrade due to bridge taps, the supplier may be forced to move a customer to a lower service tier or deny services to a customer. The customer may be unsatisfied because the system no longer meets his needs. In addition, the overall system capacity loss experienced due to the bridge taps, the moving of a customer to a lower service tier, or the denial of services to a customer, may result in financial losses for both the communications service provider and the customer.\nGiven the negative effects of bridge tap, there exists a need for mitigating the effects of bridge tap. There is also a need for methods and apparatus for minimizing the amount of time bridge tap exists during cable and/or communications device installation which may occur during, e.g., system construction, central office re-concentration, replacement projects, upgrade projects, expansion projects, and installation of back-up cables/systems to provide reserve capacity or redundancy. At least some new methods of reducing the effects of bridge tap should be suitable for use with cables which include a large number of wires commonly used in many modern applications."}
{"text": "1. Technical Field\nThis invention relates to steering systems and, more particularly, to an auxiliary steering system for vehicles for assisting a driver to laterally park a vehicle into space-limited areas.\n2. Prior Art\nConventional vehicles with wheels are classified into two categories: automobile type vehicles, whose front or rear wheels are steered to change the direction in which vehicles travel, and omni-directional vehicles, whose wheels are all steered in a certain direction so that the vehicle can travel forward, backward, right, left, or diagonally without changing the vehicle position. The conventional automobile type vehicle has a larger turning radius, which leads to the difficulty one encounters when trying to maneuver and position such a vehicle into a limited area of space. This is especially true of instances when a driver is attempting to parallel park their vehicle. During this time consuming procedure, one must often maneuver the vehicle backwards and forwards a number of times, while running the risk of striking the cars parked fore and aft of the limited space.\nSince the omni-directional vehicle can change direction without changing vehicle position or orientation, it is used in, for example, office robots, which must change direction and travel in the narrow spaces between desks. A conventional omni-directional vehicle can change direction by steering a plurality of wheels by independent steering mechanisms using special drive sources. The conventional omni-directional vehicle of this type has a steering motor for each wheel and is thus expensive. Since all the wheels must be simultaneously steered in a given direction, the steering motors must be synchronized. Synchronizing control devices are complicated and expensive.\nAccordingly, a need remains for an auxiliary steering system for vehicles in order to overcome the above noted shortcomings. The present invention satisfies such a need by providing an auxiliary steering system that is easy to use, simple in design, and eliminates the frustrations associated with parallel parking. Instead of repeatedly traveling back and forth to maneuver a vehicle into a tight parking space on a congested street, such a system enables the vehicle to be driven laterally into the parking spot. This advantageously greatly reduces the possibility of accidentally bumping into other parked cars or stationary objects. The system is also easily adaptable to a variety of vehicles and can thus be employed by many vehicle manufacturers."}
{"text": "1. Field of the Invention\nThe invention relates to a readout device, and more particularly to an array-type readout device, a dual-function readout device, and a detecting circuit for a readout device.\n2. Description of the Related Art\nIn Wan-Jun Lin, Chao P. C. P., Shir-Kuan Lin, Hsiao-Wen Zan, “A Novel Readout Circuit for an OTFD Gas Sensor with a New Front-end Trans-impedance Amplifier”, Sensors, IEEE, pp. 1141-1144, 2011, an impedance detecting circuit is proposed. However, the proposed impedance detecting circuit includes two operational amplifiers, resulting in a large size that is unfavorable for use in an array-type readout device. Such a large detecting circuit is only suitable for use in a single-type readout device, and may have a relatively low sensitivity and a relatively low signal-to-noise (SNR) ratio."}
{"text": "Temporary root canal sealers are known. Conventional temporary root canal sealers comprise a strongly alkaline aqueous mixture of calcium hydroxide and a radio-opaque filler such as barium sulfate.\nTemporary root canal sealers are intended to be a short-term means of cleaning and disinfecting the root canal cavity formed upon extraction of the dental pulp and/or root canal tissue, thereby preventing staining of the tooth upon effecting permanent sealing.\nTemporary root canal sealers containing radio-opaque materials may be employed for obtaining information about the internal environment of pulp chamber and/or associated root canals of any given tooth.\nConventional temporary root canal sealers exhibit relatively poor radio-opacities. Moreover, conventional temporary root canal sealers are poor in storage stability since the aqueous mixture of calcium hydroxide and a radio-opaque filler become heterogenous non-uniform mixtures over time.\nDE 19961002 C2 discloses a temporary root canal filling material which incorporates an X-ray contrast material, together with calcium hydroxide, an exsiccant and an organic solvent, which is an absorptive, antibacterial filler of a dry consistency, which neither hardens nor encroaches into the apical region of the root canal cavity.\nU.S. Pat. No. 5,540,766 relates to a non-aqueous dental composition comprising calcium hydroxide, gutta-percha, a radio-opaque substance and a substance which adds rigidity to the composition, which is used to make thermoplastic root canal points for obturation. Summary of the Invention\nIt is the problem of the present invention to provide a temporary root canal sealer dispersion which has improved radio-opacity, resistance against bacterial growth and consistency.\nIt is a further problem of the present invention to provide a temporary root canal sealer dispersion which is stable in both the oral environment and during storage.\nIt is a yet further problem of the present invention to provide a temporary root canal sealer dispersion which is non-toxic, anti-inflammatory and acts as a sterilant."}
{"text": "The field of the present invention relates to data mining techniques and, more particularly, to techniques for incorporating human interaction in an effective way so as to design similarity functions and perform class supervision of data.\nThe design of data mining applications has received much attention in recent years. Examples of such applications include similarity determination and classification. In the context of data mining, it is assumed that we are dealing with a data set containing N objects in a dimensionality of d. Thus, in this data space, each object X can be represented by the d coordinates (x(1), . . . x(d)). These d coordinates are also referred to as the features in the data. This is also referred to as the feature space which may reveal interesting characteristics of the data.\nThe effective design of distance functions used in similarity determination has been viewed as an important task in many data mining applications. The concept of similarity has been widely discussed in the data mining literature. A significant amount of research has been applied to similarity techniques such as, for example, those discussed in the literature: A. Hinneburg et al., xe2x80x9cWhat is the nearest neighbor in High Dimensional Space?,xe2x80x9d VLDB Conference, 2000; C. C. Aggarwal, xe2x80x9cRe-designing distance functions and distance based applications for high dimensional data,xe2x80x9d ACM SIGMOD Record, March 2001; and C. C. Aggarwal et al., xe2x80x9cReversing the dimensionality curse for similarity indexing in high dimensional space,xe2x80x9d ACM SIGKDD Conference, 2001, the disclosures of which are incorporated by reference herein.\nA different but related problem in data mining is the prediction of particular class labels from the feature attributes. In this problem, there is a set of features, and a special variable called the class variable. The class variable typically draws its value out of a discrete set of classes C(1), . . . C(k). A test instance is defined to be a data example for which only the feature variables are known, but the class variable is unknown. Training data is used in order to construct a model which relates the features in the training data to the class variable. This model can then be used in order to predict the class behavior of individual test instances, also referred to as class labeling. The problem of classification has been widely studied in the literature, e.g., J. Gehrke et al., xe2x80x9cBOAT: Optimistic Decision Tree Construction,xe2x80x9d ACM SIGMOD Conference Proceedings, pp. 169-180, 1999; J. Gehrke et al., xe2x80x9cRainForest: A Framework for Fast Decision Tree Construction of Large Data Sets,xe2x80x9d VLDB Conference Proceedings, 1998; R. Rastogi et al., xe2x80x9cPUBLIC: A Decision Tree Classifier that Integrates Building and Pruning,xe2x80x9d VLDB Conference, 1998; J. Shafer et al., xe2x80x9cSPRINT: A Scalable Parallel Classifier for Data Mining,xe2x80x9d VLDB Conference, 1996; and M. Mehta et al., xe2x80x9cSLIQ: A Fast Scalable Classifier for Data Mining,xe2x80x9d EDBT Conference, 1996, the disclosures of which are incorporated by reference herein.\nHowever, as sophisticated and, in some cases, complex as these similarity and classification techniques may be, these conventional automated techniques lack benefits that may be derived from human interaction during their design and application stages. Therefore, techniques are needed that effectively employ human interaction in order to design and/or perform data mining applications such as similarity determination and classification.\nThe present invention provides techniques for incorporating human or user interaction in accordance with the design and/or performance of data mining applications such as similarity determination and classification. Such user-centered techniques permit the mining of interesting characteristics of data in a data or feature space. For example, such interesting characteristics that may be determined in accordance with the user-centered mining techniques of the invention may include a determination of similarity among different data objects, as well as the determination of individual class labels. These techniques allow effective data mining applications to be performed in accordance with high dimensional data.\nIn accordance with a first aspect of the present invention, a computer-based technique of computing a similarity function from a data set of objects comprises the following steps/operations. First, a training set of objects is obtained. The user may preferably provide such training data. Next, the user is presented with one or more subsets of objects based on the training set of objects, wherein each subset comprises at least two objects of the data set. Preferably, the subset is a pair of objects from the data set. The user then provides feedback regarding similarity between the one or more subsets of objects. One or more sets of feature variables are defined based on features in the one or more subsets of objects. Next, one or more class variables are created in accordance with the user-provided feedback. Lastly, a similarity function or model is constructed which relates the one or more sets of feature variables to the one or more class variables.\nThus, advantageously, similarity between objects is represented as some function or algorithm determined by the attributes of the objects. The similarity model is then effectively estimated from the data set and user reactions.\nIn accordance with a second aspect of the present invention, a computer-based technique of classifying a test instance in accordance with a data set comprises the following steps. First, a test instance is obtained. The user may preferably provide such test instance. Next, the user is presented with at least one projection representing a distribution of the data set. The user then isolates a portion of the data presented in the at least one projection based on a relationship between the test instance and the data presented in the at least one projection. For instance, the user may isolate a subset of the data in the projection which the user determines to be most closely related to the test instance. Next, the behavior of the isolated portion of data is determined. Then, a class is determined for the test instance based on the isolated portion of data, when the user makes a decision to do so based on the determined behavior of the isolated portion of data. Alternatively, when the user makes a decision not to have a class determined for the test instance based on the isolated portion of data, other portions of the data set or a subset of the isolated portion of the data may be considered.\nFurther, in a preferred embodiment, the user is presented with two or more projections respectively representing different distributions of the data set such that the user may select one of the projections to be used when isolating a portion of data whose behavior is to be considered.\nThus, advantageously, such a class labeling methodology according to the invention provides a technique of decision path construction, in which the user is provided with the exploratory ability to construct a sequence of hierarchically chosen decision predicates. This technique provides a clear understanding of the classification characteristics of a given test instance. At a given node on the decision path, the user is provided with a visual or textual representation of the data in a small number of sub-spaces. This can be used in order to explore particular branches, backtrack or zoom-in into particular sub-space-specific data localities which are highly indicative of the behavior of that test instance. This process continues until the user is able to construct a path with successive zoom-ins which is sufficiently indicative of a particular class. The process of zooming-in is done with the use of visual aids, and can isolate data localities of arbitrary shapes in a given sub-space.\nIt is to be appreciated that the classification techniques of the present invention are more powerful than any of the conventional classification methods, since the invention uses a combination of computational power and human intuition so as to maximize user understanding of the classification without sacrificing discriminatory power. The result is a technique which, in most cases, can classify a test instance with a small amount of user exploration.\nThese and other objects, features and advantages of the present invention will become apparent from the following detailed description of illustrative embodiments thereof, which is to be read in connection with the accompanying drawings."}
{"text": "A method and a device for operating an internal combustion engine are discussed in DE 10 2005 051 701 A1, in which an overall injection is subdivided into a basic injection and at least one measured injection. The injection period of the measured injection is successively reduced and the injection period of the basic injection is increased, in such a way that an overall injection quantity determined from a characteristics curve of the valve remains unchanged. A deviation of a variable characterizing an actual mixture, provoked by the successive reduction of the injection period of the measured injection, from a variable characterizing a setpoint mixture is detected. The deviation or a characteristics curve of the fuel injector is adapted or corrected. The detection as to whether the actual fuel-air mixture is deviating from the setpoint fuel-air mixture take place with the aid of a Lambda value provided by a Lambda sensor."}
{"text": "Consumers continually pressure integrated circuit manufacturers to provide devices that are smaller and faster, so that more operations can be performed in a given amount of time, using fewer devices that occupy a reduced amount of space and generate less heat. For many years, the integrated circuit fabrication industry has been able to provide smaller and faster devices, which tend to double in capacity every eighteen months or so.\nHowever, as integrated circuits become smaller, the challenges of fabricating the devices tend to become greater. Fabrication processes and device configurations that didn't present any problems at a larger device size tend to resolve into new problems to be overcome as the device size is reduced. For example, in the past there was very little incentive to planarize the various layers from which integrated circuits are fabricated, and which are formed one on top of another. Because the devices themselves were relatively wide, the relatively thin layers that were formed did not present many challenges to overcome in regard to surface topography.\nHowever, as the devices have been reduced in size they have become relatively narrower. Although layer thickness has also generally decreased, the surface topography of an underlying layer tends to create greater problems for the proper formation of the overlying layer to be formed, unless the underlying layer is planarized in some way prior to the formation of the overlying layer.\nThere are several different methods used for planarizing a layer on an integrated circuit. For example, chemical mechanical polishing can be used to physically and chemically erode the surface of the layer against a polishing pad in a slurry that contains both physically and chemically abrasive materials. Further, electropolishing can be used to thin an electrically conductive layer. Unfortunately, neither process tends to produce surface topographies that are as flat as desired.\nFor example, although each of these two planarization processes tends to preferentially remove higher portions of a layer, they also attack to at least some degree the lower portions of the layer. Thus, even the though the higher portions of the layer are removed at a rate that is somewhat greater than that of the lower portions, and hence some planarization does occur, there also tends to be some amount of dishing in the lower portions of the layer, where a greater amount of material is removed than is desired.\nWhat is needed, therefore, is a method whereby the dishing of planarized layers is reduced."}
{"text": "JP2000-120810 (See FIGS. 1 and 2) discloses a conventional motor-incorporated hypocycloid-type speed reducer, in which an eccentric portion (carrier), which rotates an external gear (a planetary gear) in an eccentric manner, is connected to a rotor, which is electrically rotated, in such a manner that a rotational force of the rotor can be transmitted. The rotor is fitted with, and is supported by, a bearing rod so as to freely rotate.\nIn the above-described motor-incorporated hypocycloid-type speed reducer, the external gear is supportedly mounted on a ball bearing, and an output shaft fixed with an internal gear is rotatably supported by a plain bearing formed at a central portion of a case. According to this type of structure, when the output shaft is subjected with a load applied in a radial direction, the external gear or the internal gear (the output shaft) may on occasions tilt or lean due to the ball bearing or a clearance between the plain bearing and the output shaft, which may lead to an occurrence of a noise. Such tilting or lean of the external gear, or, of the internal gear (the output shaft) can be prevented by increasing an axial length of the plain bearing axially supporting the output shaft. This, however, may result in an increase in an axial length of the hypocycloid-type speed reducer.\nA need thus exists to provide to a motor-incorporated hypocycloid-type speed reducer, which can restrain an angle amount of tiling or lean of an output shaft."}
{"text": "The present invention relates to electrical connectors and, more particularly, electrical connectors which provide an assembler with sensory indication that a good connection has been made and which help prevent unintentional disconnection.\nU.S. Pat. No. 176,069 is exemplary of an early type of connector which was designed to provide the user with indication that a good connection has been made, and to prevent inadvertent disconnection. The '069 connector includes a female receptacle and a male plug. The receptacle has a pair of outwardly extending arms which are received within an annular groove provided on the exterior of the male plug. A set screw secures a conductor, which projects from the male plug, to the receptacle. U.S. Pat. No. 4,781,611 discloses another type of connector having a pair of deformable arms to releasably secure a female receptacle to a male spade.\nU.S. Pat. No. 2,579,739 discloses a connector having a male spade and a female receptacle. The male spade provides a pair of notched recesses, and the female receptacle has a body with a pair of deformable wings which are bent back over the body to form a pair of resilient jaws. When the spade is inserted between the body portion and the jaws, a rounded tip provided by each of the jaws extends into one of the notched recesses and connects the spade to the receptacle. Related spade to receptacle connectors are shown in U.S. Pat. Nos. 4,220,388; 4,556,747; 4,558,913; 4,720,273; and 5,038,199."}
{"text": "The present invention relates to an apparatus and a method for cutting a plastic optical fiber adapted for optical communication, which fiber has on its outer periphery a coating including a cladding layer.\nCutting of ends of plastic optical fibers for optical communication in order to mount the plastic optical fibers to optical signal connector plugs has been performed, for example, by pressing a sharp cutter to a plastic optical fiber, thereby shearing it with the pressing force (see Japanese Utility Model Registration No. 2573619). Such forcible cutting, however, causes an inconvenience that chips and/or cracks may occur in a cut plane of the plastic optical fiber, to degrade an accuracy of work in the subsequent step of forming an end face of the plastic optical fiber. To cope with such an inconvenience, an improved method has been disclosed, for example, in Japanese Patent Laid-open No. Hei 7-294748, wherein a plastic optical fiber is cut in a state that temperatures of both a cutter and the plastic optical fiber are raised by heating, thereby preventing occurrence of chips and/or cracks in a cut plane of the plastic optical fiber. Such a method, however, presents another problem that the raised temperature may give rise to deformation of a portion, other than an end face, of the plastic optical fiber, degradation of optical characteristics of the interior of the plastic optical fiber, and the like.\nIn most of plastic optical fibers (POFs), thermoplastic polymethyl methacrylate (PMMA) is used as a material of a core of the POF, and a fluorine based resin is used as a material of a cladding layer formed on an outer peripheral portion of the core. In addition, at present, a plurality of kinds of plastic optical fibers, each of which is of a multi-mode type, are being commercially available.\nIn the case of connecting a plastic optical cable led from one equipment to another equipment by mounting a POF of the cable to a connector plug and inserting the connector plug in a receptacle provided on another equipment, it is required that no change in optical transmission characteristics may occur even if the cable be replaced with a new cable. FIG. 1A is a typical sectional view showing a connector plug inserted in a receptacle. As shown in this figure, an end portion, inserted in a center hole of a connector plug 11, of a POF 1 is formed in a spherical plane R having a specific radius, and a distance xe2x80x9caxe2x80x9d between the end portion of the POF 1 and a light receiving element 11c in equipment is set to a specific length.\nThe structure shown in FIG. 1A is designed such that a positional relationship between the receptacle 11b and the connector plug 11 can be kept constant even if the connector plug 11 is replaced with a new connector plug. As a result, the replacement of the connector plug 11 with a new connector plug does not exert any effect on transmission characteristics of the POF 1 insofar as a position of the spherical plane R of the end of the POF 1 to the connector plug 11 is kept constant.\nThe formation of an end face of the POF 1 as shown in FIG. 1B is performed by making use of thermoplasticity of the POF 1. More specifically, a forming die 11d having at its one end a concave plane of a specific radius is heated and is pressed to an end face of the POF 1 mounted to the connector plug 11. The end face of the POF 1 is softened by heat of the forming die 11d, with a result that the concave plane of the forming die 11d is transferred to the end face of the POF 1 as a convex plane R.\nA volume of the softened resin is regarded not to be changed after the forming work. Accordingly, to keep constant the distance xe2x80x9caxe2x80x9d shown in FIG. 1A, the POF 1 must be mounted to the connector plug 11 such that a position of the end face to the connector plug 11 before formation of the convex plane R of the POF 1 is kept constant. For this reason, there has been adopted a method of mounting the POF 1 to the connector plug 11 with a sufficient excess portion projecting from the connector plug 11, and cutting the POF 1 at a specific position associated with the connector plug 11, thereby keeping constant the position of the convex plane R formed in the subsequent step.\nCutting of POFs has been often performed by using sharp cutting tools such as a commercially available cutter or razor. In this case, however, as shown in FIG. 1C, chips of a resin forming the POF main body and/or cracks may occur in a cut plane perpendicular to the axial line of the POF. Such chips and/or cracks may remain in an end face to be formed in the subsequent step as shown in FIG. 1B, to cause failures that change transmission characteristics, such as deficiency of a convex plane or cracking.\nIf a POF is cut with a previously heated cutter, as shown in FIG. 1D, a cladding layer and the like may extend longer in the form of whiskers, and such whiskers may adhere on a cut plane, to be rolled in a convex plane at the time of forming an end face of the POF, to degrade transmission characteristics of the POF.\nBy the way, a method of cutting a plastic material by using a cutter blade mounted to an ultrasonic vibration exciter has been known, for example, from Japanese Patent Laid-open No. Sho 58-175630, wherein the plastic material is cut by concentrating ultrasonic vibration at a portion to be cut of the plastic material via the cutter blade and imparting a force to the cutter blade being in contact with the plastic material.\nTo realize cutting of a POF with desirable optical transmission characteristics thereof kept, an attempt has been made to cut POFs by making use of such a high-frequency mechanical vibration cutting method or ultrasonic vibration cutting method. Even in the case of adopting this method, however, there may occur an inconvenience. Since a cladding layer and the like are provided on an outer peripheral surface of a core of a POF, when a cutter blade and the optical fiber are excessively heated (such heating occurs even by using a high-frequency mechanical vibration exciter), the cladding layer and the like may often extend longer in the form of whiskers and such whiskers may adhere on a cut plane, to be entrained in an end face of the POF at the time of forming the end face, thereby degrading transmission characteristics of the POF. Further, according to the related art high-frequency mechanical vibration cutting method, it is difficult to optimize working conditions required for obtaining a uniform finish cut plane, for example, a condition of determining the length of a cutoff piece of a POF.\nAn object of the present invention is to provide a method and an apparatus for cutting a plastic optical fiber, which are capable of reducing a load applied the plastic optical fiber upon cutting, thereby reducing degradation of transmission characteristics of the plastic optical fiber after cutting.\nTo achieve the above object, according to a first aspect of the present invention, there is provided a cutting apparatus for cutting a plastic optical fiber, including cutting means for cutting a plastic optical fiber, and positioning means for positioning the cutting means to at least a first position at which the plastic optical fiber is to be subjected to rough cutting, and to at least a second position at which the plastic optical fiber is to be subjected to finish cutting.\nThe cutting apparatus preferably includes vibration means for giving high-frequency mechanical vibration to a portion, being in contact with the plastic optical fiber, of the cutting means, and allowing the cutting means to heat and cut the contact portion of the plastic optical fiber.\nTo achieve the above object, according to a second aspect of the present invention, there is provided a cutting method of cutting a plastic optical fiber by using cutting means, including the steps of positioning the cutting means to a first position at which the plastic optical fiber is to be subjected to rough cutting, cutting the plastic optical fiber at the first position by the cutting means, positioning the plastic optical fiber to a second position at which the plastic optical fiber is to be subjected to finish cutting, and cutting the plastic optical fiber at the second position by the cutting means.\nIn the cutting method, preferably, in each of the cutting steps, the cutting means is heated by high-frequency mechanical vibration, to cut the plastic optical fiber.\nThe apparatus and method for cutting a plastic optical fiber according to the present invention, which are configured as described above, has the following effects:\nSince a portion to be cut of a plastic optical fiber is finely cut while being heated by a sharp cutter blade representative of the cutting means, it is possible to significantly reduce a force applied to the plastic optical fiber upon cutting, and hence to prevent occurrence of chips and/or cracks in a cut plane of the plastic optical fiber and to prolong a service life of the cutter blade and reduce a cost required for replacement and adjustment of the cutter blade. Since the plastic optical fiber is thinly cut in the finish cutting step, it is possible to form an accurate, smooth cut plane of the plastic optical fiber, and hence to improve, in the subsequent step, an accuracy of forming an end face of the plastic optical fiber while preventing occurrence of a failure in this forming step.\nIn addition to the above-described basic configurations of the present invention, according to a preferable form of the present invention, an optimum temperature control may be performed in the finish cutting step. With this configuration, it is possible to prevent occurrence of an inconvenience caused by extension of a cladding layer and the like of the softened plastic optical fiber, without the need of any additional post treatment.\nAccording to another preferable form of the present invention, the cutting apparatus having a large flexibility in working condition may be provided. More specifically, the working condition of the cutting apparatus can be changed variously in a wide range, for example, from a working condition adapted for cutting of plastic optical fibers in a very small quantity of a lot at a service shop to a working condition adapted for cutting of plastic optical fibers in a very large quantity of a lot at a factory, by changing a combination of basic cutting operations such as movement of a cutter blade representative of the cutting means in the cutting direction, movement of a plastic optical fiber in the direction perpendicular to the movement direction of the cutter blade, and control of a heating condition of the cutter blade due to high-frequency mechanical vibration.\nAccording to a further preferable form of the present invention, the cutting means may be configured as a two-piece cutter blade element having a rough cutter blade and a finish cutter blade, with the two-piece cutter blade element mounted to one cutter haft. With this configuration, it is possible to simplify the cutting apparatus and the cutting works, and hence to reduce the cost required for cutting a plastic optical fiber.\nAccording to still a further preferable form of the present invention, a plastic optical fiber may be mounted to a connector plug with a relatively longer excess portion projecting from the connector plug and the excess portion of the plastic optical fiber be repeatedly cut by several times. With this configuration, since the excess portion is repeatedly cut, it is possible to obtain a desirable cut plane of the plastic optical fiber, and since initial cutting of the excess portion of the plastic optical fiber may be roughly performed, it is possible to reduce the working cost, although the cutting is repeated by several times.\nAccording to an additional preferable form of the present invention, a peripheral device such as a cooling device may be combined with the cutting device. With this configuration, the cutting device can be easily modified into that having a function satisfying a user\"\"s specification. Further, since the cutting device can be modified into that having a function satisfying the minimum user\"\"s specification, it can be matched with each user\"\"s specification at a low cost."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device having an SOI (silicon-on-insulator) structure with a silicon layer formed on an insulator layer, and to a method of manufacturing it.\n2. Description of Related Art\nFIG. 8 shows a semiconductor device having a conventional SOI structure. As illustrated therein, an insulator layer 102 is formed on a semiconductor substrate 101, and a silicon layer 103 is formed on the insulator layer 102 to give an SOI structure.\nIn the semiconductor element region comprising the silicon layer 103, formed are source/drain regions 103a through doping with an impurity such as phosphorus, arsenic or the like, or boron or the like, and, in the area between the source/drain regions 103a, formed is a gate electrode 107 on the silicon layer 103 via a gate oxide film 106 therebetween to give a MOSFET.\nA trench 104 is formed at the element-isolation area of the silicon layer 103, and an insulating film 105 is formed within the inner wall of the trench 104 to give an element-isolation region, serving to isolate semiconductor elements from each other. As in FIG. 8, the bottom face of the silicon layer 103 makes an acute angle with the side of the element-isolation region (insulating film 105) adjacent thereto.\nAn interlayer insulating film 108 is formed on the SOI substrate having the MOSFET thereon, and a conductive layer 109 is formed on the interlayer insulating film 108. Contact holes for enabling electric connection between the conductive layer 109 and the source/drain regions 103a formed in the silicon layer 103 are formed through the interlayer insulating film 108 and filled with a conductor, via which the conductive layer 109 is electrically connected with the source/drain regions 103a. \nIn the element-isolation region of the semiconductor device having the SOI structure of that type, the bottom face of the silicon layer 103 makes an acute angle with the side of the element-isolation region (insulating film 105) adjacent thereto. In that condition, therefore, when the volume of the element-isolation region is varied through heat treatment to be effected after the formation of the element-isolation region, for example, through annealing to be effected after the formation of the oxide film 105 in the trench 104, or through heat treatment to be effected in forming the gate oxide film 106 after the formation of the element-isolation region, the volume variation shall make the silicon layer 103 have large stress at the acute-angled corners of its bottom.\nA technique for relaxing the large stress at the acute-angled corners of the bottom of the silicon layer 103 is disclosed, for example, in Unexamined Japanese Patent Publication No.(HEI)6-216230. FIG. 9 shows an SOI structure for a semiconductor device illustrated in Unexamined Japanese Patent Publication No.(HEI)6-216230. As illustrated, the trench-shaped insulator of constructing an element-isolation region is so formed that its width in the cross section is continuously increased in the downward direction, in order that the bottom of the silicon layer does not make an angle with the element-isolation region adjacent thereto. Therefore, being different from that of FIG. 8, the semiconductor device of FIG. 9 has no angle that may produce large stress, and the bottom of the silicon layer in FIG. 9 is prevented from having any large stress. In FIG. 9, numeral references are the same as those in FIG. 8, provided that the element-isolation region 105 is formed of an isolating wall 105b in the trench 104 and a polysilicon layer 105a embedded therein.\nIn the semiconductor device noted above, however, the interface of the silicon layer adjacent to the insulator is formed to be convex toward the insulator. In this, therefore, when the stress resulting from the volume variation in the element-isolation region runs toward the silicon layer, it is concentrated in some parts in the silicon layer, as will be mentioned below. The problem caused by the stress concentration is that some lattice defects are formed in those parts with the stress concentrated.\nThe reason for the stress concentration is described with reference to FIG. 10. As illustrated, in the semiconductor device of FIG. 10, the interface between the silicon layer and the insulator (trench) is so formed that, reaching the insulating layer, it is curved toward the silicon layer but not toward the element-isolation region. In this, therefore, when the stress resulting from the volume variation in the element-isolation region runs toward the silicon layer, as indicated by the arrows in FIG. 10, a plurality of stress components running in that direction shall be concentrated in the part as designated by xe2x80x9cXxe2x80x9d therein. As a result, some lattice defects are formed in that part of the silicon layer.\nMeanwhile a semiconductor device having a trench-shaped insulator of which the width is continuously decreased in the downward direction, is disclosed in IEDM (International Electron Devices Meeting) Technical Digest, 1997, p.591. However, the semiconductor device disclosed has a large depression at the surface of the trench-shaped insulator near the semiconductor element region adjacent thereto. Thus, in the case where a gate electrode extends over the trench-shaped insulator as well as the semiconductor element region, a portion of the gate electrode will be embedded in the depression of the insulator so that the distance between the semiconductor element region and the portion of the gate electrode located on the insulator will be shortened in comparison with the case of no depression. With this structure, when a transistor controllable with such a gate electrode is formed in the semiconductor element region, it will likely be influenced by an electric field from the portion of the gate electrode embedded in the depression of the adjacent insulator, that is, an electric field will be concentrated in the semiconductor element region near the depression. As a result, a reverse narrow channel effect decreasing a threshold voltage occurs and a parasitic MOSFET tends to be generated in the semiconductor element region near the depression. The concentration of an electric field may also cause a deterioration of the semiconductor element region, such as silicidation of a contact portion of the source/drain region of the transistor.\nOn the other side, Unexamined Japanese Patent Publication No. (HEI)9-8118 discloses a process for forming a trench-shaped insulator without a depression at the surface near the semiconductor element region adjacent thereto, but the width of the trench-shaped insulator formed is constant in the downward direction.\nThe object of the invention is to provide a semiconductor device in which the stress from the volume variation in the element-isolation region to the silicon layer is relaxed to thereby protect the silicon layer from having lattice defects therein and in which the surface of the trench-shaped insulator in the element-isolation region does not have a depression near the semiconductor element region, and to provide a method of manufacturing it.\nThe semiconductor device of the invention comprises an SOI substrate with a silicon layer formed on an insulating layer, and a semiconductor element region and an element-isolation region formed in the silicon layer, wherein; the element-isolation region is of a trench-shaped insulator formed adjacent to the semiconductor element region, and the trench-shaped insulator is provided with a portion which is adjacent to the semiconductor element region, of which the width is continuously decreased in the downward direction, and of which the surface is planarized near the semiconductor element region. In the device, even when the volume variation in the element-isolation region produces some stress running toward the silicon layer, the stress to the silicon layer could be well relaxed. In addition, because the surface of the trench-shaped insulator in the element-isolation region is made flat near the semiconductor element region, the concentration of an electric field in the semiconductor element region does not occur when a gate electrode is extended over the trench-shaped insulator and the semiconductor element region.\nOne method of manufacturing the semiconductor device of the invention comprises, forming, on an SOI substrate having a silicon layer formed on an insulating layer, a protective film having an unprotected portion of a predetermined pattern, forming an insulating film over the SOI substrate with the protective film, followed by etching the insulating film to leave an insulating side wall at the side of the unprotected portion of the protective film, anisotropically etching the silicon layer via the protective film and the insulating side wall both acting as a mask to form a trench in the silicon layer, and subjecting the silicon layer with the trench to a heat treatment to oxidize the side portion of the trench to thereby form a trench-shaped insulator having a width continuously decreased in the downward direction. In the method, the curved profile at the interface is controlled by varying the condition for oxidizing the side part of silicon.\nRegarding the temperature for the heat treatment in the method, if it is too low, suitable oxidation could not be attained at such low temperatures. On the contrary, if it is too high, the oxide film could be formed even at such high temperatures. In this case, however, the interface could not be curved but is linear. For these reasons, it is desirable that the temperature for the heat treatment falls between 750xc2x0 C. and 900xc2x0 C.\nAnother method of manufacturing the semiconductor device of the invention comprises, forming, on an SOI structure having a silicon layer formed on an insulating layer, a protective film having an unprotected portion of a predetermined pattern, forming an insulating film over the SOT substrate with the protective film, followed by etching the insulating film to leave an insulating side wall at the side of the unprotected portion of the protective film, anisotropically etching the silicon layer via the protective film and the insulating side wall both acting as a mask to form a trench in the silicon layer, and isotropically etching the side portion of the trench to partially remove a portion of the silicon layer underlying the insulating side wall to obtain a trench having a width continuously decreased in the downward direction, and forming an insulator in the trench. In the method, the curved profile at the interface is controlled by varying the etching condition."}
{"text": "Liquid crystalline polymer (LCP) films have properties that are very is desirable, such as excellent chemical resistance, high strength, and excellent gas barrier properties. However, these same films have certain undesirable properties. They often have poor transverse mechanical properties (i.e. they are strong in the machine direction, but tear easily in the direction transverse to the machine direction). It is also difficult to write or print on the films. LCP films are more expensive than conventional polymer films, such as polyester films.\nIt would be desirable to make multilayer films having LCP film bonded to one or more other films to obtain a film having the best properties of all of the various layers, such as a multilayer film having good gas barrier properties and relatively low cost.\nHowever, LCP films do not bond well to each other or to other films by use of an adhesive. Their surfaces do not in general adhere to adhesives. There are thus very few examples of multilayer films containing one or more LCP layers. One example can be found in Japanese Patent Application No. 02-253,950, published in 1990, where a poly(butylene terephthalate) film layer is bound to VECTRA.RTM. A 900 LCP film using a glycidyl-modified ethylene vinyl acetate adhesive. The other side of the VECTRA film is bound to polypropylene film by using two adhesive layers, a glycidyl-modified ethylene vinyl acetate layer in contact with the LCP and an ethyl acrylate-ethylene-maleic anhydride copolymer in contact with the polypropylene. Japanese Patent publications 02-253,951 (1990) and 04-135,750 (1992) use similar adhesives for binding an LCP to poly(butylene terephthalate) and polypropylene. The latter of these patent applications also uses a saponified ethylene-vinyl acetate copolymer as an adhesive. Other publications that discuss multilayer films comprising an LCP barrier layer include Japanese Patent Publication 02-307,751 (1990), PCT Patent Publication WO 95/23180, and European Patent Application No. 763,423."}
{"text": "Hearing aids are often advantageously constructed to be as small as practical given the requirements of the hearing aid. Relatively smaller hearing aids can be less intrusive to a user and less visible to others. However, as with small objects generally, relatively small hearing aids can naturally be more difficult to locate than relatively larger hearing aids and may be more easily misplaced. Additionally, as many hearing aids can be utilized in settings in proximity of many other similar hearing aids, such as clinical settings, finding one particular hearing aid that has been misplaced may be made increasingly challenging as the size of hearing aids shrink."}
{"text": "Components on spacecraft frequently are mounted on flexible isolator devices (damper struts), like the one in U.S. Pat. No. 6,003,849, configured in a hexapod arrangement in U.S. Pat. No. 5,305,981, or some other strut/truss support system. However, when a spacecraft is launched the components are sometimes restrained on a launch retention mechanism for many reasons, including decreasing deflections of the isolation systems beyond design limits. The launch retention mechanism can be an integral part of the isolation struts or a separate supporting structure. Current launch restraints pull the isolator into a fixed latch position away from the isolator\"\"s neutral in-orbit position, known as isolator bias. The bias can present difficulties when precise alignment between the isolator mounted payload and the base structure, to which the isolator is attached, is important for pre-launch preparations. Alignment adjustments between the payload and spacecraft are made on earth, where gravity sags the isolator away from the zero-gravity orientation that it will assume in outer space. This makes accurate payload positioning prior to launch problematic.\nA device is included in the isolator strut for locking the isolator in either a zero bias position (the damper is not loaded at either of its two extreme possible positions) or known/predetermined bias position and electronically overcome to unlock the strut. Another mechanical element is included to temporarily unlock the isolator without disturbing the other devices.\nObjects, benefits and features of the invention will apparent to one of ordinary skill in the art from the drawing and following description."}
{"text": "Within the last 40-50 years, the use of biocides (herbicides and pesticides) has increased dramatically. Among these biocides are the chlorinated derivatives of phenoxyacetic acid (PAA). 2,4-Dichlorophenoxyacetic acid (2,4-D), introduced in 1944, was the first phenoxy herbicide and is possibly the most widely used herbicide (W. Evans et al., J. Biochem. 1971, 122, 543-551). Although 2,4-D is not highly toxic, cleanup of inadvertently spilled 2,4-D is still necessary (P. Amy et al., Appl. Environ. Microbiol. 1985, 49, 1237-1245; R. Beadle and A. Smith, Eur. J. Biochem. 1982, 123, 323-332; W. Evans et al., J. Biochem. 1971, 122, 543-551; P. Fisher et al., Bacteriol. 1978, 135, 793-804; L. Geer et al., Appl. Environ. Microbiol. 1992, 58, 1027-1030).\nAlthough large amounts of 2,4-D are used each year, this herbicide usually does not accumulate in soil or water. Of the chlorinated phenoxy herbicides, it is less recalcitrant than more highly chlorinated compounds in the same herbicide class. 2,4-D is degraded naturally in soil under favorable environmental conditions by indigenous, competent microbial communities. There have been isolated from soil and water a number of organisms that are capable of degrading chlorinated phenoxy herbicides, such as 2,4-D, including a number of bacterial strains in multiple genera (L. Kozyreva et al., Mikrobiologiya. 1992, 62, 110-119; G. Chaudhry and G. Huang, J. Bacteriol. 1988, 170, 3897-3902; P. Amy et al., App. Environ. Microbiol. 1985, 49, 1237-1245; R. Don and J. Pemberton, J. Bacteriol. 1981, 145, 681-686; P. Fisher et al., J. Bacteriol. 1978, 135, 798-804; J. Tiedje et al., J. Agr. Food Chem. 1969, 17, 1021-1026; see also, Reineke and Knackmuss, Ann. Rev. Microbiol. 1988, 42, 263-287) and fungi (P. Donnelly et al., App. Environ. Microbiol. 1993, 59, 2642-2647\nThe best understood biochemical pathway for the microbial degradation of 2,4-D is that of Alcaligenes eutrophus JMP134 (R. Don and J. Pemberton, J. Bacteriol. 1981, 145, 681-686; R. Don et al., J. Bacteriol. 1985, 161, 85-90; R. Don and J. Pemberton, J. Bacteriol. 1985, 161, 466-468). In this pathway, 2,4-D is first converted by an .alpha.-ketoglutarate-dependent dioxygenase (the product of the tfdA gene) to 2,4-dichlorophenol and then by a DCP hydroxylase (the product of the tfdB gene) to 3,5-dichlorocathechol (DCCAT). After several additional enzymatic steps, chloromaleylacetic acid is finally produced (R. Don et al., J. Bacteriol. 1985, 161, 85-90; B. Kaphmmer and R. Olsen, J. Bacteriol. 1990, 172, 5856-5862). The genes encoding the enzyme responsible for the catabolism of 2,4-D by A. eutrophus, tfdA, tfdB and tfdCDEF, have been located on plasmid pJP4 (B. Kaphmmer and R. Olsen, J. Bacteriol. 1990, 172, 5856-5862; B. Kaphmmer et al., J. Bacteriol. 1990, 172, 2280-2286; E. Perkins and P. Lurquin, J. Bacteriol. 1988, 170, 5669-5672; W. Streber et al., J. Bacteriol. 1987, 169, 2950-2955). The expression of tfdB is regulated by the product of gene tfdS using 2,4-D and DCP as effectors (B. Kaphmmer and R. Olsen, J. Bacteriol. 1990, 172, 5856-5862). Since both 2,4-D and DCP can induce the expression of tfdB, DCP is further transformed by microorganisms bearing pJP4 or its derivatives, pRO101 or pRO103 (A. Harker et al., J. Bacteriol. 1989, 171, 314-320). No DCP can be detected by the color reaction in the culture media of these microorganisms when 2,4-D is supplied.\nAnalytical support for 2,4-D cleanup techniques usually involves gas-chromatography or high pressure liquid chromatography (HPLC) (see, e.g., K. Short et al., App. Environ. Microbiol. 1991, 57, 412-418) or immunoassays (e.g., the enzyme-linked immunosorbent assay [ELISA], 2,4-D RaPID Assay.RTM., Ohmicron, Newtown, Pa.). Even the presently available immunoassays are relatively expensive. Therefore, it is important to develop quick, inexpensive, and easy-to-use 2,4-D detection methods for use in the field, particularly methods which can be employed by non-specialists. It would be particularly advantageous to develop an assay method which could be applied beyond 2,4-D to other related chemical compounds, particularly other phenoxy biocides, many of which are less easily degraded by naturally occurring organisms."}
{"text": "In machining of metals and other materials, there are many applications where it is desirable to make a hole and thread it with a single tool. The single tool for hole making and threading is needed not only for through holes but also for blind holes in a workpiece. The advantage of such a tool is that it reduces the number of motions or operations required to produce a threaded hole; it is especially advantageous in conjunction with high speed machining operations.\nIn the prior art, it is known to use a single tool for drilling a hole and forming internal threads by swaging or coining the metal. A tool of this kind is disclosed in the Barth U.S. Pat. No. 2,703,419 granted Mar. 8, 1955, which describes a swaging tap having a drill or end mill on the leading end thereof to cut the hole which is to be threaded. This tool is rotated and the hole is cut and threaded on the inward feed of-the tool and the tool is unscrewed from the threaded hole. A similar tool is disclosed in the Grenell U.S. Pat. No. 4,271,554 granted Jun. 9, 1981. In tools of this kind, the maximum diameter of the swaging tap must be larger than the maximum diameter of the drill of end mill. Further, as described in the Grenell patent, the minor diameter of the female threads is smaller than the diameter of the drilled hole and the major diameter of the threads is larger than the diameter of the drilled hole.\nA general object of this invention is to provide a combined hole making and threading tool which may be used for both through holes and blind holes and which lends itself to high speed machining."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication system and method for selecting an HPLMN (Home Public Land Mobile Network), and more particularly to a communication system and method capable of reducing power consumption when an MS (Mobile Station) searches for an HPLMN or higher-priority PLMN.\n2. Description of the Related Art\nRoaming is defined as a function of enabling an MS (Mobile Station) to receive a service outside of a predetermined area and to receive a service anywhere within a country as a single area. A PLMN (Public Land Mobile Network) is a public telecommunication network including the MS using telephony services from mobile communication carriers through the MS's roaming function, a BTS (Base Transceiver Subsystem), a subscriber location register and an MSC (Mobile Switching Center). An HPLMN (Home PLMN) is a PLMN used at a time when the MS selects another PLMN to receive a service from a corresponding service carrier.\nIn a second or third generation mobile communication service suggested by ETSI (European Telecommunications Standards Institute) or 3GPP (Third Generation Partnership Project), an MS is recommended to search for an HPLMN according to a set time period where the MS receives the service from a VPLMN (Visitor PLMN) rather than the HPLMN. Thus, the MS searches for the HPLMN using priority information when the MS selects the PLMN in an automatic mode or manual mode on the basis of the set time period stored in an SIM (Subscriber Identity Module).\nHowever, where the MS receiving a service from the VPLMN automatically searches for a PLMN in the automatic mode, the MS located in the VPLMN must search for a higher-priority PLMN than an HPLMN or current PLMN on the basis of a set time period stored in the SIM irrespective of service conditions at a current location. That is, the MS located in the VPLMN must perform an operation of searching for another PLMN every “T” minutes. At this time, “T” as a value of a period used when the MS searches for the PLMN is a value arbitrarily set and fixed without considering service situations of the PLMN in which the MS is located.\nConventionally, the MS is in synchronization with a frequency not being monitored other than a currently monitored frequency in order to search for the PLMN. At this time, the MS detects system information contained in a synchronous frequency signal after the frequency synchronization and detects the PLMN's ID from the detected system information. While performing the above-described procedure, the MS consumes additional power. Moreover, where the MS does not utilize the service situations in the PLMN in which it is located and the set “T” value is smaller than an arbitrary reference value, the MS frequently performs a search operation for the PLMN on the basis of the set “T” value. Accordingly, there is a problem in that an available time of a charged battery for supplying the MS with power is greatly reduced.\nWhere the MS moves to an HPLMN or higher-priority PLMN during a search procedure for the PLMN, the power of the charged battery is consumed, but there is an advantage in that the MS can select a desired PLMN. However, where the MS is not located in the HPLMN or higher-priority PLMN, there are other problems in that the MS consumes unnecessary power and hence an available time of a charged battery is reduced."}
{"text": "Quantum computing refers to the field of research related to computation systems that use quantum mechanical phenomena to manipulate data. These quantum mechanical phenomena, such as superposition (in which a quantum variable can simultaneously exist in multiple different states) and entanglement (in which multiple quantum variables have related states irrespective of the distance between them in space or time), do not have analogs in the world of classical computing, and thus cannot be implemented with classical computing devices.\nQuantum computers use so-called quantum bits, referred to as qubits (both terms “bits” and “qubits” often interchangeably refer to the values that they hold as well as to the actual devices that store the values). Similar to a bit of a classical computer, at any given time, a qubit can be either 0 or 1. However, in contrast to a bit of a classical computer, a qubit can also be 0 and 1 at the same time, which is a result of superposition of quantum states—a uniquely quantum-mechanical phenomenon. Entanglement also contributes to the unique nature of qubits in that input data to a quantum processor can be spread out among entangled qubits, allowing manipulation of that data to be spread out as well: providing input data to one qubit results in that data being shared to other qubits with which the first qubit is entangled.\nCompared to well-established and thoroughly researched classical computers, quantum computing is still in its infancy, with the highest number of qubits in a solid-state quantum processor currently being below 100. One of the main challenges resides in protecting qubits from decoherence so that they can stay in their information-holding states long enough to perform the necessary calculations and read out the results. Another challenge resides in coming up with fabrication techniques that provide sufficient control of the dimensions and composition of various elements in quantum circuits."}
{"text": "1. Field of the Invention\nThe present invention is directed to a display system. More particularly, the present invention is directed to a switch plate picture frame assembly adapted to be mounted to a wall switch for the purpose of covering the wall switch and for providing a display proximate and adjacent the wall switch. The present invention further directs itself to a switch plate picture frame assembly that facilitates convenient removal and replacement of an indicia bearing member and a protective cover member retained by the assembly.\n2. Prior Art\nWall switch and outlet covers capable of displaying a variety of photographic or decorative indicia when assembled and mounted in place are well known in the prior art. U.S. Pat. No. 4,425,725, directed to a combination switch plate and photograph holder, exemplifies the prior art display systems wherein replacement of either the indicia bearing member, such as a photograph or the like, or replacement of the protective cover therefor, requires the partial disassembly of the structural elements that comprise the overall display system. For instance, in the '725 Patent, replacement of photograph 14 requires the prior removal of specially adapted photograph holder plastic cover 18 which is frictionally secured to base 20 of the display system. Further, should plastic cover 18 become damaged, scratched or marred through use, it must be replaced with an equivalent, specially adapted plastic cover which may or may not be commercially available at the time replacement thereof is required. Both of these foregoing situations present inconvenience to the user of the displayed device. These inconveniences and disadvantages are similarly attendant to the display systems described in U.S. Pat. Nos. 5,675,125, 3,953,933, 2,515,820 and 5,212,899.\nDesign Pats. 308,814 and 307,538 depict decorative or informative display members removably positioned over wall mounted electrical outlets. Since the depicted display members cover the outlet, use of the outlet requires re-positioning or removal of the display member so that the outlet can be readily accessed. Further, neither of the patents teaches a protective covering for the removable display member which is thus subjected to facial damage thereof."}
{"text": "Water supply systems have long been used in homes and buildings. A water supply system uses piping to supply fixtures such as tubs, sinks, and shower, with clean, potable hot and cold water. To do so water supply systems use pipes, water heaters, valves, outlets, storage tanks, possibly one or more pumps and other devices such as splitters to distribute water from a water inlet to the various fixtures. A well designed water supply system can reliably supply clean water to users at designed rates and temperatures.\nWhile modern water supply systems have proven themselves to be highly valuable, and while they are very often required in homes and other buildings, water supply systems are not without problems. For example, some parts of the United States and many parts of the world face water shortages so water cannot be wasted. In addition, energy, water heating costs can be of major concerns in some locations and applications. These are problems because to obtain hot water from a fixture can require running cold or tepid water from the hot water fixture until the water line is cleared of cooler water and the hot water line fills with hot water.\nThe cold or tepid water usually just goes down a drain. While a minor problem if the hot water fixture is near the water heater, if the fixture is a distance away significant amounts of water can be wasted. Not only is the running water lost its residual heat is also lost. So a hot water fixture fifty feet (50 ft.) away from a hot water heater may cause the loss of a significant amount of water and residual heat until hot water arrives at a suitable temperature.\nThe foregoing problems may be intolerable, particularly as water shortages continue to worsen. Therefore, a need exists for advanced water supply systems that drain water from a hot water fixture into a water recovery storage tank before water comes from the hot water fixture, that then re-fills the hot water line with hot water to supply the hot water fixture with water at the proper temperature, and then returns the drained water to a hot water heater."}
{"text": "A recent vehicle includes various electronic apparatuses mounted for convenience of passengers within the vehicle. For example, various electronic devices mounted in a vehicle include a navigation system, a multimedia system, and an air conditioning system.\nThe existing electronic devices inside the vehicle provide user interfaces through designated buttons, and a touch screen is generally used in recent days. However, based on the mounting arrangement of the electronic devices, the driver needs to check and operate the electronic devices with his/her eyes while driving, thereby reducing safe driving.\nAccordingly, there is a need for a development of an interface technology for minimizing an eye control of electronic devices of a vehicle by a driver so as to enable the driver to operate the electronic device inside the vehicle without reducing driving safety and improve a driving convenience.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the present disclosure and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "The present invention relates, in general, to the field of precision scalar network analyzers. More particularly, the present invention relates to a swept microwave power measurement system and method which determines \"absolute\" power by applying an interpolated correction factor determined at any given frequency during a sweep to the \"apparent\" power output at a predetermined frequency.\nIn the past, scalar analyzers have been used to measure magnitude only reflection and transmission characteristics of microwave components. The relatively low cost, high speed and simple operation of scaler analyzers have been their fundamental attraction. However, their limited power measurement accuracy (+/-1 to 2 dB) have often kept them in the role of qualitative limit line testing or intermediate screening devices instead of being capable of more precise quantitative analysis.\nTherefore, it has been necessary to use scalar analyzers in association with accurate, high performance power meters when accurate power measurement is required. This is necessitated by virtue of the fact that a conventional scalar network analyzer does not have the capability of accurately measuring microwave power, since scalar measurements are, per se, power ratio measurements. It is for this reason that active component manufacturers are forced to make power measurements with a separate instrument such as a microwave power meter. In addition to the added equipment cost, this operation involves considerable extra measurement time per part. Nevertheless, the combination of a scalar analyzer and a power meter still renders it very difficult to make exact measurements of power level at many different frequencies without laborious effort on the part of an operator and swept absolute measurements are impossible.\nIt would therefore be highly desirable to incorporate an accurate power measurement capability into a scalar analyzer. Further, the ability to make high performance, swept absolute power measurements would be a fundamental advantage. Absolute power measurements are required to set up critical drive levels to components such as mixers, or to measure critical output levels such as the compression point of amplifiers. And, while absolutely flat linearity is essential for ratio and relative measurements conventionally made by scalar analyzers, more is demanded for absolute power measurements. The two major specifications for such a system would be an accurate, traceable 1 mW calibration reference to set absolute power at a known frequency (which is typically 50Mhz) as well as an accurate frequency response curve for the power sensor."}
{"text": "Identifying the flows generated by different application-layer protocols is of major interest for network operators. Such identification enables QoS (quality of service) engineering for different types of traffic, such as voice and video traffic, and enables specific applications such as traffic forensics applications, network security applications, etc. Moreover, it enables ISPs to control resource intensive applications, such as peer-to-peer (P2P) applications, to limit and/or control application traffic and usage. Similarly, in enterprise networks, it is very important for administrators to know activities on their network, such as services that users are running, the applications dominating network traffic, etc.\nThroughout this disclosure, the term “flow” refers to a sequence of packets exchanged between two network nodes, referred to as a source and a destination of the flow where the source or the destination may be the originator of the exchange. Generally, in an IP network, such as the Internet, a flow is identified by a 5-tuple of <source IP address, destination IP address, source port, destination port, protocol> where the payload of the flow may be represented by a string of alphanumeric characters and other sequences of bits."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an automobile anti-theft apparatus. More specifically the present invention relates to a gas projection apparatus for generating a train of gas ring vortexes which may be aimed at a moving vehicle to stall the vehicle.\n2. Description of the Prior Art\nThe theft of automobiles, trucks and like moving vehicles is one of the most serious crimes in the United States today. The loss to the economy of such automobile thefts is estimated to be in the hundreds of millions of dollar per year. Insurance premiums, already high, have skyrocketed as a result of such theft.\nIn addition, a new phenomena known as \"car jacking\" in which an automobile is stolen while an innocent victim, who is driving the automobile, is, for example, at a stop light or is parking the automobile. Car jackings often result in injury or even the death of the victim of the car jacking.\nIn the past a number of devices have been developed to assist in the prevention of thefts of automobiles, trucks and the like. These devices may be generally classified as active devices or passive devices. An active device is one which disables a system in the automobile, normally required to drive the vehicle. There are active systems which disable steering columns, fuel systems, transmissions and drive train systems, as well as steering wheels themselves.\nFor example, one such prior art active anti-theft device disclosed in U.S. Pat. No. 4,790,406 comprises a stainless shield permanently fixed to the shaft bowl of the steering column of an automobile. When assembled on a steering column, the shield moves the shaft bowl when the gear shift is moved. A lock assembly with a deadbolt is provided on the defender shield with this deadbolt mounted for radial movement with respect to the shield. The lock is positioned such that when the car is in the parked position the deadbolt enters the gate of the shaft device at a position such that the gear lever cannot be moved out of the park position.\nA second active vehicle anti-theft system disclosed in U.S. Pat. No. 4,762,198 comprises an auxiliary switch that is interposed in the ignition controlling circuit between the main key switch and the ignition starting unit. The auxiliary switch is formed with normally open single pole, single break contacts located behind and spaced from the outer surface of a vehicle wall and access to the contacts is through a single small diameter entrance tube, which functions also to mount the switch behind the vehicle wall. The switch is provided with a separate probe which is insertable in the entrance tube so as to pass through the vehicle wall into engagement with the contacts of the switch, acting to close the contacts.\nThese and other prior art active anti-theft automobile systems function well to prevent theft of the automobile when the automobile is parked. However when the automobile is at a stop light and a car jacking occurs, the active anti-theft automobile systems currently available are unable to prevent a theft of the automobile and the possible harm to the driver of the automobile.\nOther vehicle anti-theft systems are of a passive nature and do not disable the vehicle. These include audio alarms which alert anyone within listening range that there has been an unauthorized tampering with the lock systems or otherwise unauthorized entry into the vehicle. However, audio alarms are generally easily disabled by the thief. In addition, such alarm systems are frequently ignored as they are often inadvertently activated, such as by innocent passers-by.\nWhat is needed is a relatively simple, yet highly effective vehicle anti-theft system which will effectively disable a moving vehicle thereby preventing vehicle theft and car jackings."}
{"text": "The prevalence of overweight people in US has reached alarming levels. Also the proportion of children and adolescents who are overweight has tripled in the past three decades.\nObesity arises as a consequence of positive caloric balance. A comprehensive behavioral approach comprising a gradual increase of energy expenditure from exercise and an appropriate diet to decrease the caloric intake should be the more effective treatment of obesity.\nHowever, this approach has a relatively low success rate. Consequently alternative forms of treatment, including surgery and/or medication, have been developed in an effort to increase the likelihood of achieving, and maintaining weight loss. In particular pharmacotherapy, in combination with intensive behavioral treatment, can lead to clinically significant decreases in body weight in obese population.\nThe FDA-approved weight-loss drugs are phentermine, sibutramine, orlistat and diethylpropion. Among them, phentermine is one of most efficient and safe in promoting weight loss especially when given along with recommendations for diet.\nPhentermine is a sympathomimetic amine which first received approval from the FDA in 1959 as an appetite suppressant for the short-term treatment of exogenous obesity for patients with an initial body mass index ≧30 kg/m2, or ≧27 kg/m2 in the presence of other risk factors (e.g., hypertension, diabetes, hyperlipidemia).\nPhentermine hydrochloride ({circumflex over (α)}, {circumflex over (α)}-dimethylphenethylamine hydrochloride) became available in the United States in the early seventies and is currently sold in several dosage form such as tablets, film coated tablets and capsules. Orally disintegrating tablets (ODT) which dissolve in the mouth for oral or sublingual administration are dosage forms particularly useful for patients with swallowing problems, for example children.\nIn particular, buccal tablets are intended for disintegrating in the mouth; the patient places them in the buccal cavity on the tongue or between cheeks and gums, thereby allowing a slow dissolution, which usually require 30-60 minutes (E. Rotteglia: “Compresse farmaceutiche” Societá Editoriale Farmaceutica, Milan, Italy, 1966).\nOn the contrary, sublingual tablets are intended to be placed under the tongue, where the active ingredient can be directly absorbed through the mucosa. These forms are provided with slow-disintegrating formulation as well (E. Rotteglia, ibid. and S. Casadio, Technologia Farmaceutica II Ed., Cisalpina Goliardica, Milan, Italy).\nOrally disintegrating tablets with this kind of prolonged release are hardly suitable for formulating active ingredients, such as analgesics or anti inflammatory agents, which have to exert an immediate effect. Also, they are not always suitable for patients such as children or elderly people, and for the administration of active ingredients with an unpleasant taste because of the long stay in the mouth.\nSublingual tablets with a rapid dissolution profile can be prepared according to the Zydis® freeze-drying procedure. Zydis® is a registered trademark of R. P. Scherer Company (Manufacturing Chemist, February 1990). However, such formulations are very expensive and require sophisticated technologies and methods from the production point of view. These products are substantially freeze-dried products, the pharmaceutical formulation being therefore difficult to handle (due to its friability and fragility) and requiring specific packaging. A problem with freeze-dried sublingual tablet formulations is the impossibility to effect any taste-masking on the active ingredient.\nWO088/08298 (Fuisz Technologies) discloses rapid-dissolution pharmaceutical composition in which the active ingredient is included in a water-soluble carrier obtained through a specific preparation process which requires a specific, expensive plant. Moreover, the resulting compositions exhibit friability problems and must always be handled and packed with particular precautions (use of dehydrating agents, humidity-tight packages, controlled-humidity work environmental and so on).\nEP-A-494972 (Cima Labs Inc.) describes effervescent tablets suitable to the direct oral administration, i.e. without a previous development of the effervescence in water, consisting of microcapsules containing the active ingredients and an amount of effervescent agent sufficient to promote the release of microgranules when ingested and to give a “fizzing” sensation when in contact with the buccal mucosa of the patient.\nIn this case also, notwithstanding the presence of amounts of effervescent agent lower than in conventional formulations (60% by weight compared with the total composition) the typical cautions used for effervescent tablets should be taken.\nFurther drawbacks are the friability of the tablets and the use of microcapsules. In fact, the preparation technique described in EP-A-494972 does not foresee any wet granulation, i.e. using a solvent, but the direct mixing of the powder and its subsequent compression. Such a preparation technique yields tablets having friability values higher than those involving wet granulation of the mixture to be compressed."}
{"text": "The term crosstalk was originally coined to indicate the presence in a telephone receiver of unwanted speech sounds from another telephone conversation. Of particular interest is crosstalk that is caused by signal coupling between adjacent circuits. The most common coupling is due to near-field effects and can usually be characterized by mutual inductance and direct capacitance. This is best illustrated by considering two parallel balanced transmission paths. One circuit (the disturbing circuit) is a source of signal energy that is undesirably coupled into an adjacent circuit via stray capacitance and mutual inductance. Near-end crosstalk NEXT) is crosstalk energy that travels in the opposite direction to that of the signal in the disturbing circuit, whereas far-end crosstalk is crosstalk energy that travels in the same direction as the signal in the disturbing circuit. Circuit analysis indicates that NEXT is frequency dependent and, for communication connectors, its magnitude typically increases with frequency at a 6.0 dB per octave rate NEXT is introduced within an electrical cable as a result of signal energy being coupled between nearby wires; and within an electrical connector, particularly modular plugs and jacks, as a result of signal energy being coupled between nearby conductors. NEXT is undesirable and is frequently referred to as offending crosstalk.\nU.S. Pat. No. 5,096,442 discloses a modular jack whose NEXT is about 25 dB below the level of the incoming signal at 100 MHz. Such NEXT is attributable to crosstalk that is introduced by the combination of a standard modular plug with a standard modular jack such as are generally used for voice-grade communications. However, this level of crosstalk is generally too high for modern high-speed data applications.\nU.S. Pat. No. 5,186,647 discloses a substantial improvement to the design of a standard modular jack by crossing the path of one of the conductors within the jack, over the path of another of the conductors within the jack to produce crosstalk of an opposite polarity. Such compensating crosstalk attempts to cancel NEXT rather than merely minimizing it by, for example, increasing the separation between conductors. This simple technique improves NEXT at 100 MHz by a startling 17 dB, thereby enabling popular modular jacks to meet Category 5 requirements specified in ANSI/EIA/TIA-568A. An example of such a modular jack is the M100 Communication Outlet, which is manufactured by Lucent Technologies Inc.\nTechniques have been developed that further improve the crosstalk performance of an electrical connector so that NEXT is now more than 60 dB below the level of the incoming signal at 100 MHz. U.S. Pat. No. 5,997,358 shows such techniques. However, this level of crosstalk performance represents the very best that can be attained since crosstalk will vary according to how the plug is seated within the jack. At least one manufacturer has disposed the jack springs within the modular jack at a relatively large contact angle (about 36.degree.) with respect to the longitudinal axis of the modular jack in order to push the modular plug into a fixed location within the jack. However, since there are many jack springs that need to make electrical contact with the blades of an inserted modular plug, large contact angles make this task difficult. Whereas large contact angles create increased pressure against the plug blades, increased pressure by some of the jack springs can preclude other spring contacts from making contact with the plug blades unless the plug blades and the jack springs are all precisely aligned. Indeed, current FCC standards recommend a relatively small contact angle (i.e., between 13 and 24 degrees) to assure that all plug blades make contact with the jack springs.\nAccordingly, what the prior art appears to lack and what is now desired is a technique for assuring the consistent positioning of a modular plug within a modular jack, where the modular jack includes jack springs that are disposed at relatively small angles with respect to the longitudinal axis of the jack."}
{"text": "1. Cross-Reference to Related Applications\nThis application is a related to U.S. patent application Ser. No. 07/917,199, filed Jul. 22, 1992, entitled: \"METHOD FOR RETORTING ORGANIC MATTER\", which is a CIP of U.S. patent application Ser. No. 07/820,134, filed Jan. 13, 1992, entitled: \"METHOD FOR RETORTING ORGANIC MATTER\", and is also related to U.S. patent application Ser. No. 07/917,191, filed Jul. 22, 1992, entitled: \"APPARATUS FOR RETORTING ORGANIC MATTER\".\n2. Field of the Invention\nThe present invention relates generally to an apparatus for allowing thermal dimensional changes of metal parts in a retort mechanism, and more particularly to an apparatus which allows radial and axial dimensional changes of a retort auger and associated equipment during the heating and cooling stages of a retort mechanism while maintaining a sealed relationship with a retort chamber.\n3. Description of Prior Art\nIn the past, there have been many methods and apparatus for disposing of or treating waste materials and contaminated materials for recycling. Procedures have been utilized in the past for cleaning up contaminated materials and for recycling materials containing hydrocarbons. Such prior methods have included chemically treating the materials, burning the materials, disposing of the materials in landfills, and retorting the materials under high temperatures. Some examples of prior methods and apparatus as disclosed in the prior art are described as follows.\nU.S. Pat. No. 3,682,115 to Rodgers, issued Aug. 8, 1972, discloses a portable disposal apparatus for combustible waste in which the combustible waste is crushed and chopped and conveyed to a combustion chamber where it is ignited with an auxiliary fuel and burned. Products of combustion which have not been fully consumed are condensed in condenser tanks. Unburned gases are then directed back into the combustion chamber to sustain combustion while residual tars, oils, and condensed liquids are removed from the condenser tanks from time to time.\nU.S. Pat. No. 4,235,676 to Chambers, issued Nov. 25, 1980, discloses an apparatus including an elongated tube that is maintained at a temperature of about 1100 degrees Fahrenheit and through which organic waste material, such as shredded rubber automobile tires or industrial plastic waste or residential trash which preferably has metal and inorganic matter removed therefrom, is moved at a uniform rate of speed in the absence of air and/or oxygen. The vapors and gases which are produced and/or liberated within the tube are quickly removed therefrom by means of a vacuum of from about four inches to about six inches of mercury, with the vapors being condensed and the gases separated therefrom.\nU.S. Pat. No. 4,308,103 to Rotter, issued Dec. 29, 1981, discloses a system including a cylindrical, horizontally disposed reactor vessel having a material conveying device including a plurality of paddle-like impellers mounted on a rotatable pipe for transporting comminuted solid carbonizable materials, such as coal, shredded scrap tires, comminuted municipal waste, sawdust and wood shavings, and the like, through the reactor vessel; a heating chamber arranged coaxially around the reactor vessel to subject the material passing through the reactor vessel to an indirect heat transfer relationship with a burning air-fuel mixture spirally swirling within the heating chamber and moving in a direction generally countercurrent to the material passing through the reaction vessel with the burning air-fuel mixture and combusted gases being progressively constricted and confined by the heating chamber. One end of the reaction vessel has a feed material inlet. Communicating with the feed material inlet is a gravity packed feed material column which assists in effectively sealing the feed material inlet from oxygen-containing gases. A rotary air lock is located near the upper end of the feed material column to further assure the exclusion of oxygen-containing gases from the interior of the reaction vessel. A side inlet may be provided in the feed material column for introducing inert gas to such column to further seal the system against oxygen-containing gases. The other end of the reaction vessel has a solid residue outlet and a gas-vapor outlet. Communicating with the solid residue outlet is a gravity packed column which contributes to the sealing at the outlet end of the reaction vessel from oxygen-containing gases. The comminuted solid carbonizable material passing through the reactor vessel is converted by pyrolysis, or high-temperature destructive distillation, into combustible gases, liquid hydrocarbons and solid carbonaceous residues. Gases and vaporized liquids generated from the solid carbonizable material introduced into the reaction vessel leave the reaction vessel through the gas-vapor outlet and are withdrawn under a slightly negative pressure and in a manner so as to avoid the entrance of any oxygen-containing gases into the reaction vessel.\nU.S. Pat. No. 4,715,965 to Sigerson et al., issued Dec. 29, 1987, discloses a method for separating volatilizable contaminants from soil by introducing the soil into a rotary aggregate dryer through which a working gas indirectly heated to between 750 and 1800 degrees Fahrenheit is drawn to vaporize the contaminants, and for recovering the contaminants for disposal or for cooling the effluent to condense and precipitate out a substantial portion of the contaminants and passing the effluent through activated carbon.\nU.S. Pat. No. 4,730,564 to Abboud, issued Mar. 15, 1988, discloses a multi-stage rotary kiln for burning waste and including a pair of concentric tubes affixed one inside the other with waste being conveyed through the inner tube and with hot burning gases being introduced into the inner tube to cause the waste to burn.\nU S. Pat. No. 4,821,653 to Jones, issued Apr. 18, 1989, discloses an apparatus for detoxifying heavy metals and the like contained in sludges, soils, incinerated ashes and similar materials by passing the metal-containing material through a pyrolyzer means operated with a substantially oxygen-free environment.\nU.S. Pat. No. 4,974,528 to Barcell, issued Dec. 4, 1990, discloses a method for removing hydrocarbon contaminants from soil by advancing the soil through a dryer having a combustion chamber therein, and exposing the soil to a gaseous flame in the combustion chamber to volatilize certain of the contaminants in the soil.\nThe patents to Rodgers, Abboud and Barcell teach direct contact between a flame and the material being treated.\nThe patents to Chambers and Jones teach an anaerobic treatment.\nThe patent to Sigerson et al. teaches drawing a hot working gas stream at a temperature of between 750 degrees Fahrenheit and 1800 degrees Fahrenheit through the soil by an induced draft fan.\nThe patent to Rotter teaches moving a spiralling high temperature heating medium within the heating zone toward the material inlet end of the reaction pipe."}
{"text": "Almost all physiological processes are based on molecular recognition of peptides or proteins and other biologically active components and the like. A lot of peptides having important biological functions such as hormones, enzymes, inhibitors, enzyme substrates, neurotransmitters, immunomodulators and the like have been found to date. There are resultantly many studies conducted to develop therapeutic means with a peptide, with understanding physiological effects of active substances composed of these peptides.\nIn development of a peptide as a medicinal product, there are new methods established for treatments and therapies of diseases correlated with peptides, however, in use of a peptide as a medicinal product, problems as described below are generated. That is, a) under physiological conditions, most peptides are decomposed by specific and nonspecific peptidases, to give low metabolic stability, b) due to relatively large molecular weight thereof, absorption after ingestion is poor, c) excretion through liver and kidney is fast, and d) since a peptide is structurally flexible and receptors for a peptide can be distributed widely in an organism, undesired side effects occur in non-targeted tissues and organs.\nExcept for some examples, relatively small natural peptides (peptide composed of 30 to less than 50 amino acids) are present under disorderly conditions due to a lot of conformations in dynamic equilibrium in a diluted aqueous solution, as a result, the peptides lack in selectivity for a receptor and become liable to undergo metabolism, thus, determination of a biologically active conformation is difficult. When a peptide itself has a biologically active conformation, namely, when having the same conformation as that under condition linked to a receptor, a reduction in entropy in linking to a receptor is smaller as compared with a flexible peptide, consequently, an increase in affinity to a receptor is expected. Therefore, there is a need for a biologically active peptide having a uniformly controlled conformation, and development thereof is important.\nThere are recently many efforts conducted to develop a peptide mimic or a peptide analog (hereinafter, referred to as “peptide mimic” together) showing a more preferable pharmacological property than that of a natural peptide as the original form thereof. “Peptide mimic” used in the present specification is a compound which is capable of mimicking (agonistic substance) or blocking (antagonistic substance) at receptor level, the biological effect of a peptide, as a ligand of a receptor. For obtaining a peptide mimic as the most possible agonistic substance, factors such as a) metabolic stability, b) excellent bioavailability, c) high receptor affinity and receptor selectivity, d) minimum side effects, and the like should be taken into consideration. From the pharmacological and medical standpoint, it is often desirable not only to mimic the effect of a peptide at receptor level (agonistic action) but also, if necessary, to block a receptor (antagonistic action). The same items as the pharmacological items which should be considered for designing a peptide mimic as the above-described agonistic substance can be applied also to designing of a peptide antagonistic substance.\nOne example of peptide mimics is development of a peptide having a controlled conformation. This mimics, as correctly as possible, a conformation linked to a receptor of an endogenic peptide ligand. When analogs of these types are investigated, resistance to a protease increases, and resultantly, metabolic stability rises and selectivity rises, thereby lowering side effects.\nOverall control in the conformation of a peptide is possible by restricting flexibility of a peptide chain by cyclization. Cyclization of a biologically active peptide not only improves its metabolic stability and selectivity for a receptor but also gives a uniform conformation, thereby enabling analysis of the conformation of a peptide. The cyclization form is the same as that observed in natural cyclic peptides. Examples thereof include side chain-side chain cyclization, or side chain-end group cyclization. For cyclization, side chains of amino acids not correlated with receptor recognition can be mutually linked, or can be linked to the peptide main chain. As another embodiment, there is head to tail cyclization, and in this case, a completely cyclic peptide is obtained.\nFor these cyclization operations, a cross-linking technology is imperative. Typical examples of cyclization include cross-linkages via a disulfide bond (SS bond), an amide bond, a thioether bond and an olefin bond. More specific examples thereof include cyclization by connecting two penicillamine residues via a disulfide cross-linkage (Mosberg et al., P.N.A.S. US, 80:5871, 1983), cyclization by forming an amide bond between lysine and aspartic acid (Flora et al., Bioorg. Med. Chem. Lett. 15 (2005) 1065-1068), a procedure in which an amino acid derivative containing a cross-linked portion having a thioether bond introduced previously is introduced into a peptide bond and cyclization thereof is performed in the last condensation reaction (Melin et al., U.S. Pat. No. 6,143,722), and cyclization by mutually cross-linking (S)-α-2′-pentenylalanines introduced into the main chain using an olefin metathesis reaction (Schafmeister et al., J. Am. Chem. Soc., 122, 5891-5892, 2000).\nA cross-linkage via a disulfide bond, however, will be cleaved by a reductase generally present in an organism. Also, a cross-linkage via an amide bond will be cleaved by an enzyme cutting an amide structure present in an organism. A thioether bond and an olefin bond need substitution of side chains of an amino acid in a peptide elongation process, for attaining cyclization thereof.\nAlso known is a cross-linked structure originating from nitrogen constituting an amide in the peptide main chain skeleton, as a method needing no modification of a side chain of a peptide (Gilon et al., Biopolymers 31:745, 1991). However, this peptide will be cleaved by an enzyme cutting an amide structure, because of inclusion of an amide bond in this peptide.\nFurther, known as a cross-linked peptide having a molecular structure capable of linking to other substituent is a cross-linked peptide utilizing 2,4,6-trichloro[1,3,5]-triazine (Scharn et al., J. Org. Chem. 2001, 66, 507-513). In this method, however, the reaction in forming a cross-linked portion is an aromatic nucleophilic substitution reaction, thereby limiting applicable peptides.\nAs the analogous peptide, a cross-linked peptide in which a side chain and a carboxy terminus are linked is known (Goodman et al., J. Org. Chem. 2002, 67, 8820-8826). This peptide, however, will be cleaved by an enzyme cutting an amide structure, because of inclusion of an amide bond in this peptide.\nA peptide having a controlled conformation is expected to provide a lot of pharmacological use applications. For example, somatostatin is a cyclic tetradecapeptide present in both the central nerve system and surrounding tissues and has been identified as an important inhibitor against secretion of a grow hormone from pituitary gland, and additionally, has functions such as suppression of secretion of glucagon and insulin from spleen, regulation of most gastrointestinal hormones, regulation of release of other neurotransmitters correlated with motor activity and a recognition process all over the central nerve system, and the like. A cross-linked peptide composed of nine amino acids called a WP9QY (W9) peptide mimicking the steric structure of a contact site between TNF and a TNF receptor suppresses the inflammation activity of TNFa, and additionally, is known to suppress bone resorption (Aoki et al., J. Clin. Invest. 2006; 116(6):1525-1534).\nThere is a study conducted to obtain a peptide mimic having metabolic stability improved by adding to the peptide a structure not present in natural peptides, as the peptide mimic showing a more preferable pharmacological property than that of a natural peptide as the original form thereof, in addition to a cross-linked peptide having a conformation controlled as described above. For example, resistance to an enzyme is improved by using cross-linkages (a cross-linkage via a thioether bond, a cross-linkage via an olefin, and the like) other than the above-described natural cross-linking (disulfide cross-linkage). Further, resistance to metabolism in an organism is improved by adding, for example, PEG and the like, to the terminus or the side chain of a peptide.\nJP-A No. 2004-59509, compounds described in PCT international publication WO2007/034812, compounds described in PCT international publication WO2007/122847, compounds described in PCT international publication WO2010/104169 and compounds described in PCT international publication WO2010/113939"}
{"text": "The present invention relates to a network system in which a plurality of internetwork apparatuses such as routers each connecting networks at the network layer level are used to connect a plurality of networks to one another, and more particularly to a network system having the redundant configuration in which a current internetwork apparatus is changed to a standby internetwork apparatus upon failure of the current internetwork apparatus.\nIn the system in which a plurality of networks are connected by a single router, when a failure occurs in the router, operation of the whole of the system is stopped, so that the whole of the system fails to be operated normally. On the other hand, there is a system having the redundant configuration using two routers one serving as a current router and the other serving as a standby router and in which in the normal state of the current router the system is operated by the current router and when a failure occurs in the current router the system is operated by the standby router instead of the current router to thereby attain a high reliability system.\nSuch a system having the redundant configuration is disclosed in, for example, JP-A-6-131208. In this system, when a failure occurs in a current network apparatus, operation of the whole of the current network apparatus is stopped and the current network apparatus is changed to a standby network apparatus to thereby realize the system having the redundant configuration.\nFurther, in a network system, various management information such as statistical information, set information and the like is exchanged among nodes in the network and is managed as a management information base (MIB).\nIn the system having the redundant configuration in the prior art, when a failure occurs in a port of the current network apparatus or in a connection portion between the port and the network, operation of the whole of the current apparatus is stopped. That is, operation of not only the port in which the failure occurs but also other all normal ports in the current network apparatus is stopped and operation is changed from the current network apparatus to the standby network apparatus. Accordingly, even in communication between the networks connected to the normal ports of the current network apparatus, it is necessary to disconnect the normal ports and to change the communication routes to the standby network apparatus.\nFurther, in this case, when data is transmitted from a terminal of the network, the relay port of the data is changed and accordingly an address thereof must be also changed.\nIn addition, even in transmission and reception of management information managed by using management information base (MIB), it is necessary to transmit and receive the management information in consideration of the physical configuration of the network apparatus in response to change of the route.\nMore particularly, in the system for managing networks on the basis of MIB information, that is, information based on the management information base (MIB) for each network, when a failure occurs in a certain port of a current apparatus, operation of the current apparatus is stopped completely and is changed from the current apparatus to the standby apparatus. Accordingly, MIB information of all ports (that is, all networks) of the current apparatus cannot be obtained. Thus, management of the networks cannot be performed on the basis of the MIB information available before occurrence of the failure.\nIt is an object of the present invention to provide a network system and a failure restoration method of the network system which is adapted to be solve the drawbacks in the prior art.\nIt is another object of the present invention to provide a network system and a failure restoration method of the network system in which when a failure occurs in a port of a current apparatus, data communication can be made without change of other normal ports of the current apparatus.\nIt is a further object of the present invention to provide a network system and a failure restoration method of the network system in which even a terminal having no dynamic routing function and address resolution protocol (ARP) function for conversion of address, that is, having only the routing function can make communication without considering change of a route due to occurrence of a failure.\nIt is another object of the present invention to provide a network system capable of collecting management information base (MIB) information without considering change of a route due to occurrence of a failure.\nIn order to achieve the above objects, according to an aspect of the present invention, there is provided a network system which includes a plurality of networks, a first internetwork apparatus including a plurality of first ports each connected to the plurality of networks, a second internetwork apparatus including a plurality of second ports each connected to the plurality of networks, and a data transmission unit connected to the first and second internetwork apparatuses to transmit data mutually between the first and second internetwork apparatuses, wherein the first internetwork apparatus includes a first transmitting and receiving unit connected to the plurality of first ports and the data transmission to transmit and receive data mutually among the plurality of first ports and the data transmission unit, a failure detection unit for detecting whether a failure occurs in any of one of the plurality of first ports and a route between one of the plurality of first ports and one of the plurality of networks connected thereto and, when occurrence of a failure is detected, issuing notification indicating a failure occurrence portion to transmit the notification through the first transmitting and receiving unit to the data transmission unit, and a first port control unit for causing each of the plurality of first ports to be able to transmit and receive data to and from the one of the plurality of networks in a normal state and causing the one of the plurality of first ports not to be able to transmit and receive data to and from the one of the plurality of networks in response to detection of occurrence of failure by the failure detection unit, and the second internetwork apparatus includes a second transmitting and receiving unit connected to the plurality of second ports and the data transmission unit for transmitting and receiving data mutually among the plurality of second ports and the data transmission unit and for receiving the notification transmitted through the data transmission unit to produce the notification, and a second port control unit for causing the plurality of second ports not to be able to receive data from the plurality of networks and to be able to transmit data to the plurality of networks in the normal state and responsive to the notification from the second transmitting and receiving unit to cause one of the plurality of second ports connected through the one of the plurality of networks to the one of the plurality of first ports to be able to transmit and receive data to and from the one of the plurality of networks.\nAccording to one example of the present invention, the first port control unit sets an address of each of the plurality of first ports to a common address common to the port and a corresponding one of the plurality of second ports connected through a corresponding one of the plurality of networks in the normal state, the common address being different for each of the plurality of first ports, and the second port control unit sets an address of each of the plurality of second ports to an address inherent to the port in the normal state.\nAccording to one example of the present invention, the first port control unit includes a unit for disconnecting the one of the plurality of first ports from the one of the plurality of networks in response to detection of occurrent of failure by the failure detection unit, and the second port control unit changes an address of the one of the plurality of second ports from the inherent address to the common address set in the one of the plurality of first ports in response to the notification from the second transmitting and receiving unit.\nAs described above, even if a failure occurs in any portion of a port in the current internetwork apparatus (router), the routing module including the port and a route between the port and the pertinent networks, only the port corresponding to the portion where the failure occurs is stopped or electrically disconnected from the pertinent network and the normal port in the current system is operated as it is. Further, the address of a port of the standby system corresponding to the failed port where the failure occurs is changed to an address of the failed port, that is, the port of the standby system is caused to be able to perform transmission and reception. Thus, change of the transmission route of packet data can be minimized and communication between the networks can be continued. In other words, even if a failure occurs in a portion of the current router, it is not necessary to change the whole current system to the standby system as in the prior art, the route corresponding to only the failed portion is changed and the routes in the normal portions of the current system are not required to be changed.\nFurther, the relay address of the packet data from a terminal is not required to be changed and the route is automatically changed so that the packet data is transmitted to the terminal of the destination.\nIn other words, since the relay address of the packet data from the terminal is not required to be changed and the packet data from the terminal is automatically transmitted to the destination terminal, the terminal can perform communication without considering change of the route for the communication.\nAs described above, according to the present invention, the reliable and inexpensive LAN system can be constructed by realization of duplication at a unit of port. Further, the terminal having neither dynamic routing function nor ARP function can communicate without considering change of a route. In addition, management of the network such as collection of statistical information and the like can be made without considering the redundant configuration.\nIn order to achieve the above objects, according to another aspect of the present invention, there is provided a network system further comprising a management terminal connected to the one of the plurality of networks for managing management information base (MIB) information of each of the plurality of first and second ports, wherein the first internetwork apparatus further comprises a first MIB information memory unit for storing the management information base information for each of the plurality of first ports, and a first MIB information control unit connected to the first MIB information memory unit and the first transmitting and receiving unit, and the second internetwork apparatus further comprises a second MIB information memory unit for storing management information base (MIB) information for each of the plurality of second ports, and a second MIB information control unit connected to the second MIB information memory unit and the second transmitting and receiving unit, the management terminal sending a collection request of management information base information of the one of the plurality of first ports through the one of the plurality of networks, the one of the plurality of second ports and the second transmitting and receiving unit to the second MIB information control unit after the one of the plurality of second ports has been caused to be able to transmit and receive data to and from the one of the plurality of networks, the second MIB information control unit being responsive to the collection request of the management information base information from the management terminal to read out management information base information from the second MIB information memory unit after the one of the plurality of second ports has been caused to be able to transmit and receive data, and sending a collection request of management information base information available before occurrence of failure of the one of the plurality of first ports through the second transmitting and receiving unit, the data transmission unit and the first transmitting and receiving unit to the first MIB information control unit, the first MIB information control unit being responsive to the collection request of the management information base information from the second MIB information control unit to read out management information base information available before occurrence of failure of the one of the plurality of first ports from the first MIB information memory unit and transmitting the management information base information through the first transmitting and receiving unit, the data transmission unit and the second transmitting and receiving unit to the second MIB information control unit, the second MIB information control unit calculating a sum total of management information base information read out from the first and second MIB information memory units to send the sum total through the second transmitting and receiving unit, the one of the plurality of second ports and the one of the plurality of networks to the management terminal.\nAs described above, when the collection request of the statistical information (MIB) is issued from the management terminal of the network to the system having the redundant configuration, the internetwork apparatus operating as the current system reads out the statistical information of the internetwork apparatus serving as the standby system and the statistical information of the whole redundant configuration system is collectively returned to the management terminal. Thus, the management terminal can manage the networks without considering the physical redundant configuration.\nAs described above, even if a failure occurs in a port of the current system, operation of only the failed port (or all ports included in the routing module of the failed port) is stopped and other normal ports of the current system are continued to be operated, so that the port of the standby system corresponding to the failed port (or the stopped port) is operated. Accordingly, the past management information base (MIB) concerning the stopped port of the current system can be obtained. Hence, the networks can be managed on the basis of the management information base (MIB) information available before occurrence of the failure.\nFurther, since the address of the port where the failure occurs is used as the address of the port of the standby system operated instead, the management terminal can obtain the management information base (MIB) information of the port of the current system of which operation is automatically stopped without changing the address of the source requiring the MIB information regardless of change of the route."}
{"text": "Ethernet provides high speed data communications between two nodes that operate according to the IEEE 802 Ethernet Standard. The communications medium between the two nodes can be twisted pair wires for Ethernet, or other types of communications mediums that are appropriate. PoE systems provide power and data over a common communications link. More specifically, a power source device (PSE) coupled to the physical layer of a first node of the communications link provides direct current (DC) power (for example, 48 volts DC) to a powered device (PD) at a second node of the communications link. The DC power is transmitted simultaneously over the same communications medium with the high speed data from one node to the other node.\nExample PDs that utilize PoE include Internet Protocol (IP) phones, and wireless access points, etc. The PSE typically includes a serializer/deserializer (i.e. SERDES) coupled to a transceiver, and/or a physical-layer (PHY) device, to support high speed serial data transport. Herein, data ports and their corresponding links can be interchangeably referred to as data channels, communication links, data links, etc, for ease of discussion.\nPoE integrated circuits are sensitive to harmful surge events, such as electrostatic discharge (ESD) and cable discharge events (CDE). The PHY of a PoE system is particularly susceptible to damage due to ESD and CDE. During these surge events, currents can be extremely high (e.g., 100 amps) and it becomes vital to ensure that voltages do not exceed critical breakdown and spark gap limits of PoE systems and their respective circuits.\nAn ESD event typically occurs when a device becomes charged as a result of mishandling or improper packaging and then discharged by a sudden connection to ground. CDE, on the other hand, can occur when a charge accumulates on a cable, such as a twisted pair cable used in Ethernet networks, and is connected to an Ethernet port of lower potential. The resulting high-energy discharge may damage the device to which the cable is connected. Coupling of external events, like a lightning strike, is yet another example source of CDE.\nSome PoE integrated circuits have a conventional protection circuit to combat detrimental ESD and CDE. These conventional protection circuits discharge electrostatic or harmful surge energy using a capacitor and/or a transient-voltage-suppression (TVS) diode. Unfortunately, when managing ESD and CDE events on multiple ports, conventional PoE protection circuits require at least one capacitor per port to provide a low impedance path to ground, and/or at least one TVS diode per port to supply protection from differential and common mode transients. These capacitors are commonly referred to as bulk capacitors and are typically large and consume a significant amount of circuit board space, as well as contribute to additional cost. Similarly, the need for multiple TVS diodes contributes to circuit board space requirements and cost.\nThus, what is needed is a protection circuit for PoE devices that overcomes the shortcomings described above."}
{"text": "This invention relates to a novel calixarene derivative which is soluble in organic solvents, films of the novel compound and a method of forming a pattern in a film formed by applying a solution of the novel compound to a substrate.\nCalixarenes are cyclic oligomers formed by condensation of phenols and formaldehyde, and some calixarene derivatives can be obtained by substitution reactions after the condensation reaction.\nRecent studies have revealed that calixarenes and their derivatives, like cyclodextrine and crown ethers, have the ability to form inclusion compounds. Synthesis of water soluble calixarenes has achieved a limited success, and it is under study to use calixarene derivatives as adsorbents for the recovery of, for example, uranyl ion from seawater or heavy metal ions from waste water.\nThere is good expectation that calixarenes and their derivatives will serve as advantageous functional materials. From a practical point of view, functional materials are generally required to be soluble in ordinary organic solvents and capable of providing films from solutions. However, known calixarenes and their derivatives are very low in solubilities in organic solvent, viz. below 1 wt %, and hence it is hardly conceivable to practically use these compounds in the form of films as functional materials. J. Am. Chem. Soc., Vol. 111 (1989), 8192-8200 shows to form very thin, monolayer-like films of some calixarene derivatives by using the Langmuir-Blodgett technique, but practical applications of such films will be quite limited."}
{"text": "1. Field of Invention\nThis invention relates to a baling apparatus, and more particularly to a square baler with plunger providing increased impact force.\n2. Description of Related Art\nConventional square hay balers include a bale forming chamber and a reciprocating plunger that slides into and out of the chamber. As the chamber receives loose hay material, the plunger slides into the chamber during a compaction stroke to compress the loose hay material into the form of a bale. Such balers typically include a drive train that transmits power to the reciprocating plunger.\nHowever, prior art square hay balers have certain deficiencies. For instance, conventional hay balers apply very high forces to the plunger in order to complete the compaction stroke. Because these forces are transmitted by the drive train and by connecting rods that connect the drive train to the plunger, the connecting rods and drive train components must be designed to accommodate this loading."}
{"text": "It is well known that water supplied to households from some wells and community water supplies often tastes unpleasant or is dangerous to drink because of minerals, chemicals, organisms and organic materials that are dissolved or suspended in the water. Widespread recognition of this information accounts for the variety of domestic water purification devices that have been developed and patented. However, use of water purifiers in homes is not as widespread as one would expect in view of the scope of the problem.\nMost water distillers developed for home use have an electric heating element immersed in raw (that is, undistilled) water that is supplied to a boiler-evaporator from the water mains of a dwelling. The mass of water in the boiler is raised to boiling temperature. The resulting steam is conducted through a fin-type condenser coil from which the distillate emerges. In some designs a motor driven fan forces ambient air over the condenser fins for cooling and condensing the steam. In other designs the condenser coil is water cooled by locating it in a chamber into which the raw, comparatively cool water is fed before the raw water is conducted to the boiler, resulting in waste water. Most distillers on the market distill on a batch-by-batch basis rather than continually according to demand, as should be the case.\nAmong the reasons that installations of previously existing distiller designs have been small in number, although there is such a great need for them, is that the distillers are configured in a way that makes them difficult to install in a concealed and inconspicuous manner near the kitchen sink, where water is usually consumed in the home. Since an existing type of distiller would ordinarily be installed near the kitchen sink, one possibility is to stand the prior art types of distillers on a counter top next to the sink. Yet, most householders object to dedicating to a distiller precious counter top area, which is usually felt to be insufficient in most residences in the first place. Installation inconvenience becomes a factor in deciding not to buy any distiller presently on the market. Besides, most, if not all known prior art distillers can be characterized as lacking any redeeming aesthetic characteristics.\nAnother place in which a prior art distiller might be installed is in a cabinet near the kitchen sink. The problem with this is that prior art distillers are vertically oriented, that is, they have one component stacked on another so they have a tall profile or an inordinately great height dimension. As a practical matter, this means that they require dedication of a lot of below-the-counter top cabinet space, and it becomes impossible to use any of the space in the cabinet, above, below or on the sides of the distiller.\nBesides deficiencies in aesthetic characteristics and excessive space utilization, prior art distillers are difficult to maintain in good operating condition, particularly, because of the difficulty of cleaning sediment and scale from the internal parts of the distiller. Most prior distillers require a substantial amount of disassembly and handling or working on multiple parts to fully clean the boiler of scale. This may be an aggravating factor that the user realizes only when the distiller fails to produce distilled water up to rated capacity. Facilitating easy and simple descaling and cleaning are problems that have frequently been attacked but have not been completely solved in prior distiller designs."}
{"text": "1. Field of the Invention\nThe present invention relates to a growth factor which interacts with the human oncogene erbB-2, and which stimulates as well as inhibits the growth of cells overexpressing this oncogene. A ligand is described which is capable of binding to the expression product of the erbB-2 oncogene. The present invention additionally relates to anti-ligand molecules capable of recognizing and binding to the erbB-2 ligand molecule and to screening assays for such ligands. The present invention further relates to uses for the erbB-2 ligand, the anti-ligand molecules and the screening assay. Furthermore, the invention relates to a cloned gene capable of expressing the erbB-2 ligand of the present invention.\n2. Description of the Related Art\nCarcinogenesis is believed to be a multi-step process of alteration of genes which are involved in the growth control of cells. A variety of proto-oncogenes and oncogenes have been implicated in the activation of tumor cells as regulating factors. For example, oncogenic protein kinases are believed to induce cellular transformation through either inappropriate or excessive protein phosphorylation, resulting in the uncontrolled growth of malignant neoplasms. See Wrba, F., et al., Histopathology, 15:71-76 (1989).\nOne group of proto-oncogenes encodes cellular growth factors or their receptors. The c-erbB-1 gene encodes the epidermal growth factor or its receptors. The c-sis gene encodes the B-chain of the platelet-derived growth factor. The c-fms gene encodes a related or identical molecule for the receptor of the granulocyte-macrophage colony stimulating factor. A fourth member of this group of proto-oncogenes, called neu was identified in ethylnitrosourea-induced rat neuroblastomas.\nThe human counterpart of neu, called HER-2/neu or c-erbB-2, has been sequenced and mapped to the chromosomal locus 17q21. See Schneider, P. M., et al., Cancer Research, 49:4968-4971 (Sep. 15, 1989). The HER-2/neu or c-erbB-2 oncogene belongs to the erbB-like oncogene group, and is related to, but distinct from the epidermal growth factor receptor (EGFR). The c-erbB-2 oncogene is known to express a 185 kDa transmembrane glycoprotein (p185.sup.erbB-2). The expressed protein has been suggested to be a growth factor receptor due to its structural homology with EGFR. However, known EGFR ligands, such as EGF or TGF.alpha., do not bind to p185.sup.erbB-2.\nThe oncogene has been demonstrated to be implicated in a number of human adenocarcinomas leading to elevated levels of expression of the p185 protein product. For example, the oncogene has been found to be amplified in breast, ovarian, gastric and even lung adenocarcinomas. Furthermore, the amplification of the c-erbB-2 oncogene has been found in many cases to be a significant, if not the most significant, predictor of both overall survival time and time to relapse in patients suffering from such forms of cancer. Carcinoma of the breast and ovary account for approximately one-third of all cancers occurring in women and together are responsible for approximately one-fourth of cancer-related deaths in females. Significantly, the c-erbB-2 oncogene has been found to be amplified in 25 to 30% of human primary breast cancers. See Slamon, D., et al., Science, 244, 707-712 (May 12, 1989).\nAlthough ligands for EGFR are known, namely EGF and TGF.alpha., few ligands for the oncogene-encoding transmembrane proteins such as erbB-2, ros, etc., have been characterized. Transforming growth factor ligands belong to a family of heat and acid-stable polypeptides which allow cells to assume a transformed morphology and form progressively growing colonies in anchorage-independent growth assays (DeLarco, et al., Proc. Natl. Acad. Sci. USA, 75:4001-4005 (1978); Moses, et al., Cancer Res., 41:2842-2848 (1981); Ozanne, et al., J. Cell. Physiol., 105:163-180 (1980); Roberts, et al., Proc. Natl. Acad. Sci. USA, 77:3494-3498 (1980)). The epidermal growth factor receptor (EGFR) and its physiologic ligands, epidermal growth factor (EGF) and transforming growth factor .alpha. (TGF.alpha.), play a prominent role in the growth regulation of many normal and malignant cell types (Carpenter, G., Annu. Rev., Biochem., 56:881-914 (1987)).\nOne role the EGF receptor system may play in the oncogenic growth of cells is through autocrine-stimulated growth. If cells express the EGFR and secrete EGF and/or TGF.alpha., then such cells could stimulate their own growth. Since some human breast cancer cell lines and tumors express EGFR (Osborne, et al., J. Clin. Endo. Metab., 55:86-93 (1982); Fitzpatrick, et al., Cancer Res., 44:3442-3447 (1984); Filmus, et al., Biochem. Biophys. Res. Commun., 128:898-905 (1985); Davidson, et al., Mol. Endocrinol, 1:216-223 (1987); Sainsbury, et al., Lancet, 1:1398-1402 (1987); Perez, et al., Cancer Res. Treat., 4:189-193 (1984)) and secrete TGF.alpha. (Bates, et al., Cancer Res., 46:1707-1713 (1986); Bates, et al., Mol. Endocrinol, 2:543-555 (1988)), an autocrine growth stimulatory pathway has been proposed in breast cancer (Lippman, et al., Breast Cancer Res. Treat., 7:59-70 (1986)).\nThe erbB-2 proto-oncogene amplification has been found in breast, ovarian, gastric, salivary gland, and in non-small cell carcinomas of the lung (King, et al., Science, 229:974 (1985); Slamon, et al., Science, 244:707 (1989); Yokota, et al., Lancet, 1:765 (1986); Fukushige, et al., Mol. Cell. Biol., 6:955 (1986); Semba, et al., Proc. Natl. Acad. Sci. USA, 82:6497 (1985); Weiner, et al., Cancer Res., 50:421 (1990)). Amplification and/or overexpression of the erbB-2 proto-oncogene has been found to correlate with poor prognosis in breast, ovarian and non-small cell lung carcinomas (Slamon, et al., Science, 235:177 (1986); Slamon, et al., Science, 244:707 (1989); Guerin, et al., Oncogene Research, 3:21 (1988); Wright, et al., Cancer Res., 49:2087 (1989); Kern, et al., Cancer Res., 50:5184 (1990); DiFiore, et al., Science, 237:178 (1987)). In addition to these clinical studies, in vitro studies strongly suggest that overexpression of the erbB-2 transmembrane receptor (p185.sup.erbB-2) may have an important role in tumor progression (DiFiore, et al., Science, 237:178 (1987); Hudziak, et al., Proc. Natl. Acad. Sci. USA, 84:7159 (1987)).\nAn autocrine growth stimulatory pathway analogous with that proposed for epidermal growth factor receptor and its ligands may also be employed by a growing list of oncogene encoded transmembrane proteins that have structure reminiscent of growth factor receptors. This list includes the protooncogenes neu and its human equivalent erbB-2 or HER2 (Bargmann, et al., Nature, 319:226-229 (1986); Coussens, et al., Science, 230:1131-1139 (1985); Yamamoto, et al., Nature, 319:230-234 (1986); c-kit (Yarden, et al., EMBO, 6:341-3351 (1987); ros (Neckameyer, et al., Mol. Cell. Biol. 6:1478-1486 (1986); met (Park, et al., PNAS, 84:6379-6383 (1987); trk (Martin-Zanca, et al., Nature, 319:743-748 (1986); and ret (Takahashi, et al., Mol. Cell. Biol., 7:1378-1385 (1987)). The erbB-2 and c-kit protooncogenes encode factors that display remarkable structural homology with EGFR (Yarden, et al., Annu. Rev. Biochem., 57:443-478 (1988). Although erbB-2 and its related oncogene neu are related to EGFR, these proteins are distinct. For example, known EGFR ligands such as EGF and TGF.alpha. do not bind to erbB-2 receptor. (King, et al., EMBO, 7:1647 (1988); and Stern, et al., EMBO, 7:995 (1988).\nIf, according to the autocrine growth stimulatory pathway, malignant cells are capable of secreting a potent tumor growth factor in vivo, it is plausible that the growth factor ligand might be detected in body fluids, much like human chorionic gonadotropin or .alpha.-fetoprotein, and could be used as a tumor marker and a prognostic variable. Studies suggest that TGF.alpha. activity can be detected in body fluids of cancer patients and that its presence may provide important information concerning the biology of a patient's tumor (Stromberg, et al., J. Cell. Biochem., 32:247-259 (1986); Twardzick, et al., J. Natl. Cancer Inst., 69:793-798 (1982); Sherwin, et al., Cancer Res., 43:403-407 (1983)).\nPrior to the present invention, no ligand was known which binds to p185.sup.erbB- 2 protein. Thus, a need continues to exist for a ligand for p185.sup.erbB-2. Such a ligand might be used to counteract the effects of c-erbB-2 oncogene overexpression in facilitating carcinogenesis."}
{"text": "1. Field of the Invention\nThe present invention relates to a stereo slide viewer, and more specifically, relates to a stereo slide viewer in which the shape of a slide mount holder section has been improved, for better loading of the stereo slide mount and better positioning.\n2. Description of the Related Art\nThe present applicant has filed an application for this type of stereo slide mount and stereo slide viewer (Japanese Patent Application No. Hei 11-273315), and this prior application will be first described with reference to FIG. 4 to FIG. 9.\nIn FIG. 4, reference symbol 1 denotes a stereo slide viewer, and this stereo slide viewer 1 is provided with a pair of left and right ocular lenses 3L and 3R at the rear of a slide mount holder section 2 of a channel shape, and provided with an illumination lamp 4 and a circular reflector 5 in front of the slide mount holder section 2, so that by lighting the illumination lamp 4, light is irradiated onto the backside of the slide mount holder section 2 by the circular reflector 5. Ocular lens holders 6L and 6R for the left and right ocular lenses 3L and 3R have a construction capable of adjusting the distance between optical axes and focus in the direction of the optical axis.\nA stereo slide mount 7 shown in FIG. 5 is freely loadable into the slide mount holder section 2. For positioning of the stereo slide mount 7 in the back and forth direction, a flat spring 8 is provided on the front wall surface 2a of the slide mount holder section 2, and the stereo slide mount 7 is pressed against the rear wall surface 2b of the slide mount holder section 2.\nMoreover, a film pitch adjusting apparatus 9 is arranged in the middle of the left and right ocular lenses 3L and 3R, and outside of the visual field of the ocular lenses 3L and 3R. This film pitch adjusting apparatus 9 is for adjusting the interval between left and right film holders 11L and 11R of the stereo slide mount 7 loaded into the slide mount holder section 2, by rotating helical cams 10L and 10R of the film pitch adjusting apparatus 9. These helical cams 10L and 10R are formed substantially in a semicircular shape as seen from the side, and constructed such that in the initial position, the helical cams 10L and 10R do not come into the slide mount holder section 2.\nThe film pitch adjusting operation and film mounting procedure will be described next. Two left and right film holders 11L and 11R are mounted on a base frame 12 of the stereo slide mount 7, and the left and right film holders 11L and 11R are brought closest to each other, to thereby mount films 13L and 13R on the film holders 11L and 11R. The base frame 12 has left and right windows 14L and 14R, and the left and right film holders 11L and 11R have windows 15L and 15R, respectively, so that the light of the illumination lamp 4 reaches the films 13L and 13R.\nThen, a collimation pattern mask 16 shown in FIG. 6 is attached to the base frame 12 having the films 13L and 13R mounted thereon. The collimation pattern mask 16 is provided with a square window 17 so that the helical cams 10L and 10R can freely enter into and come out, as well as left and right transparent windows 18L and 18R. The same collimation pattern CP mainly composed of a plurality of vertical lines is respectively printed on the backside of the transparent windows 18L and 18R.\nSubsequently, the base frame 12 attached with the collimation pattern mask 16 is loaded into the slide mount holder section 2 of the stereo slide viewer 1 shown in FIG. 4.\nThen, sense of intimacy of images in the films 13L and 13R with respect to the collimation pattern CP of the collimation pattern mask 16 is observed, in a stereoscopic vision through the ocular lenses 3L and 3R of the stereo slide viewer 1. If the stereoscopic images are seen in the same plane as the collimation patterns CP or beyond the collimation patterns CP, the film pitch of the stereo slide mount 7 is appropriate, and adjustment of the film pitch is not necessary.\nIf the film pitch is not appropriate, a knob 19 of the stereo slide viewer 1 is rotated to adjust the film pitch. That is to say, if the knob 19 is rotated in the clockwise direction, the helical cams 10L and 10R rotate from the initial position to enter into the slide mount holder section 2, so that the points of the cam surfaces 10aL and 10aR enter into the two film holders 11L and 11R to push the film holders 11L and 11R outwards. As a result, the interval between the film holders 11L and 11R increases.\nAt this time, the area where the outside edges of the screen of the films 13L and 13R are masked by the edges of the transparent windows 18L and 18R of the collimation pattern mask 16 increases, and the distance of the stereoscopic image moves backwards with respect to the collimation patterns CP. Then, at the time when the stereoscopic image is seen in the same plane as the collimation patterns CP or beyond the collimation patterns CP, it is the optimum film pitch. At this time, the knob 19 is reversely rotated to thereby return the helical cams 11L and 10R to the initial position. During this time, the film holders 11L and 11R are held in the optimum film pitch position.\nSubsequently, the stereo slide mount 7 is taken out from the slide mount holder section 2, to remove the collimation pattern mask 16 from the base frame 12, and instead thereof, a cover frame 21 comprising two windows 20L and 20R shown in FIG. 7 is loaded. As a result, the stereo slide mount 7 having the optimum film pitch is completed.\nThe completed stereo slide mount 7 can be loaded into the slide mount holder section 2 of the stereo slide viewer 1 and used for appreciating the films. Even if the knob 19 is rotated at the time of appreciation, since the film holders 11L and 11R are covered with the cover frame 21, the helical cams 10L and 10R are not brought into contact with the film holders 11L and 11R, and hence the film pitch does not change.\nMoreover, at the time of adjusting the film pitch, in the case where the film pitch exceeds the optimum pitch and become excessive, and hence the stereoscopic image is seen far away from the collimation patterns CP, CP the film pitch is readjusted. At this time, the knob 19 is reversely rotated to return the helical cams 10L and 10R to the initial position. The stereo slide mount 7 is taken out from the slide mount holder section 2, and the left and right film holders 11L and 11R are made to slide to the position where the film holders 11L and 11R are in the closest position. Then, the stereo slide mount 7 is reloaded, so as to readjust the film pitch by rotating the knob 26.\nThe stereo slide viewer 1 is a stereo slide viewer equipped with the film pitch adjusting apparatus 9, but other than the stereo slide viewer 1, a relatively cheap stereo slide viewer (not shown) used mainly for appreciation, without having the film pitch adjusting apparatus, is also widely known.\nFIG. 8(a) shows the stereo slide mount 7, which can be loaded into either a stereo slide viewer 1 having the film pitch adjusting apparatus 9, or a stereo slide viewer not having the film pitch adjusting apparatus and used mainly for appreciation. In more detail, the left and right opposite ends are formed in a circular arc. If the vertical dimension of the stereo slide mount 7 is designated as H, the lateral minimum dimension is designated as Wa, the lateral maximum dimension is designated as wd, and the radius of the circular arc of the left and right opposite ends is designated as r, then, the stereo slide mount 7 is formed in such a relation that:\nWd={(Wa)2+(H)2}1/2\nr=[{(Wa)2+(H)2}1/2/2.\nOn the other hand, FIG. 8(b) shows another rectangular stereo slide mount 7xe2x80x2, which is mainly loaded into a stereo slide viewer used for appreciation, without having the film pitch adjusting apparatus. It is formed in a rectangular shape, with the vertical dimension being H, and the lateral dimension being Wb, which is the same size as the lateral minimum dimension Wa of the stereo slide mount 7.\nThe slide mount holder section 2 shown in FIG. 9 is provided in both of the stereo slide viewer 1 having the film pitch adjusting apparatus 9, and the stereo slide viewer not having the film pitch adjusting apparatus and used mainly for appreciation, and both of the stereo slide mounts 7 and 7xe2x80x2 can be loaded therein. The slide mount holder section 2 is formed such that the left and right internal walls of the slide mount holder section 2 are in parallel with each other, and the dimension Wg between the left and right internal walls is substantially the same as the lateral maximum dimension Wd of the stereo slide mount 7 having a large lateral dimension.\nTherefore, as shown in FIG. 9(a), the circular-arc stereo slide mount 7 can be smoothly loaded into the slide mount holder section 2, and positioning is easily performed in such a manner that the portion of the stereo slide mount 7 having the lateral maximum dimension comes in contact with the left and right internal walls of the slide mount holder section 2.\nHowever, as shown in FIG. 9(b), the rectangular stereo slide mount 7xe2x80x2 has a lateral dimension Wb shorter than the dimension Wg between the left and right internal walls of the slide mount holder section 2. Hence, an adapter 22 for positioning is required in the slide mount holder section 2. Accordingly, when the stereo slide mount 7xe2x80x2 is used, the cost for the adapter 22 is required additionally, and there is another problem in that when the stereo slide mount 7 is used by replacing the stereo slide mount 7xe2x80x2, the adapter 22 must be removed, and the adapter 22 must be stored.\nTherefore, there is a technical problem to be solved in the stereo slide viewer, so that a rectangular stereo slide mount and a stereo slide mount having circular-arc left and right opposite ends can be smoothly loaded and positioned in the slide mount holder section, without using any adapter. It is therefore an object of the present invention to solve this problem.\nThe present invention has been proposed in order to achieve the above-described object, and provides:\na stereo slide viewer comprising left and right ocular lenses, and a grooved slide mount holder section for selectively loading a rectangular stereo slide mount and a stereo slide mount having the left and right opposite ends in a circular-arc shape, in front of the left and right ocular lenses, wherein left and right internal walls of the slide mount holder section are formed slantwise such that the distance between the internal walls is narrowed in the loading direction of the stereo slide mount, and the distance between the left and right internal walls in the internal bottom portion of the slide mount holder section is formed so as to become substantially the same size as the left-right length of the loading point of the stereo slide mount, so that the loading point of the stereo slide mount is positioned with the left and right internal walls in the internal bottom portion of the slide mount holder section;\na stereo slide viewer comprising left and right ocular lenses, and a grooved slide mount holder section for selectively loading a rectangular stereo slide mount and a stereo slide mount having the left and right opposite ends in a circular-arc shape, in front of the left and right ocular lenses, wherein left and right internal walls of the slide mount holder section are formed such that the left and right internal walls are provided in parallel over a predetermined length in the loading direction from a loading port of the stereo slide mount or formed slantwise so as to be narrowed in the loading direction of the stereo slide mount, as well as being formed slantwise so as to be further narrowed in the loading direction from the position of the predetermined length, and the distance between the left and right internal walls in the internal bottom portion of the slide mount holder section is formed so as to become substantially the same size as the left-right length of the loading point of the stereo slide mount, so that the loading point of the rectangular stereo slide mount is positioned with the left and right internal walls in the internal bottom portion of the slide mount holder section, and the opposite ends of the circular-arc stereo slide mount is positioned with the left and right internal walls in the predetermined length; and\na stereo slide viewer comprising left and right ocular lenses, and a grooved slide mount holder section for selectively loading a rectangular stereo slide mount and a stereo slide mount having the left and right opposite ends in a circular-arc shape, in front of the left and right ocular lenses, wherein left and right internal walls of the slide mount holder section are formed in a curved shape protruding inwards over a predetermined length in the loading direction from a loading port of the stereo slide mount, and formed slantwise so as to be narrowed in the loading direction from the position of the predetermined length, and the distance between the left and right internal walls in the internal bottom portion of the slide mount holder section is formed so as to become substantially the same size as the left-right length of the loading point of the stereo slide mount, so that the loading point of the rectangular stereo slide mount is positioned with the left and right internal walls in the internal bottom portion of the slide mount holder section, and the opposite ends of the circular-arc stereo slide mount is positioned with the curved left and right internal walls in the predetermined length.\nAs will be described in detail in a first embodiment, the invention according to a first aspect is a stereo slide viewer, wherein left and right internal walls of the slide mount holder section are formed slantwise such that the distance between the internal walls is narrowed in the loading direction of the stereo slide mount, and the distance between the left and right internal walls in the internal bottom portion of the slide mount holder section is formed so as to become substantially the same size as the left-right length of the loading point of the stereo slide mount, so that the loading point of the stereo slide mount is positioned with the left and right internal walls in the internal bottom portion of the slide mount holder section. As a result, the stereo slide mount can be smoothly loaded into the slide mount holder section, and positioning thereof can be easily performed, without using an adapter in the conventional case.\nMoreover, the invention according to a second aspect is a stereo slide viewer, wherein left and right internal walls of the slide mount holder section are formed such that the left and right internal walls are provided in parallel over a predetermined length in the loading direction from a loading port of the stereo slide mount or formed slantwise so as to be narrowed in the loading direction of the stereo slide mount, as well as being formed slantwise so as to be further narrowed in the loading direction from the position of the predetermined length, and the distance between the left and right internal walls in the internal bottom portion of the slide mount holder section is formed so as to become substantially the same size as the left-right length of the loading point of the stereo slide mount, so that the loading point of the rectangular stereo slide mount is positioned with the left and right internal walls in the internal bottom portion of the slide mount holder section, and the opposite ends of the circular-arc stereo slide mount is positioned with the left and right internal walls in the predetermined length. As a result, the same effects as those of the invention according to the first aspect can be expected. Further, since the left and right internal walls and the stereo slide mount are brought into contact with each other in the predetermined length, more accurate positioning can be performed, and the positioning condition is maintained. Hence, it can be prevented that the stereo slide mount comes up from the slide mount holder section.\nFurthermore, the invention according to a third aspect is a stereo slide viewer, wherein left and right internal walls of the slide mount holder section are formed in a curved shape protruding inwards over a predetermined length in the loading direction from a loading port of the stereo slide mount, and formed slantwise so as to be narrowed in the loading direction from the position of the predetermined length, and the distance between the left and right internal walls in the internal bottom portion of the slide mount holder section is formed so as to become substantially the same size as the left-right length of the loading point of the stereo slide mount, so that the loading point of the rectangular stereo slide mount is positioned with the left and right internal walls in the internal bottom portion of the slide mount holder section, and the opposite ends of the circular-arc stereo slide mount is positioned with the curved left and right internal walls in the predetermined length. As a result, the same effects as those of the invention according to the second aspect can be expected, and hence this invention exhibits enormous effects."}
{"text": "The present invention relates to wireless communication systems. More specifically, the present invention relates to the control of contention-based wireless access in communication systems.\nFIG. 1 illustrates a simplified wireless spread spectrum code division multiple access (CDMA) or time division duplex (TDD) or frequency division duplex (FDD) communication system 18. The system 18 comprises a plurality of Node Bs 26, 32, 34, a plurality of radio network controllers (RNCs) 36, 38, 40, a plurality of UEs 20, 22, 24 and a core network 46. The plurality of Node Bs are connected to the plurality of RNCs 36, 38, 40, which are, in turn, connected to the core network 46. Each Node B 26, 32, 34 communicates with its associated user equipment (UE) 20, 22, 24. Data signals are communicated between UEs and the Node B over the same spread spectrum. Each data signal in the shared spectrum is spread with a unique chip code sequence. Upon reception, using a replica of the chip code sequence, a particular data signal is recovered.\nIn the context of a CDMA system, signals are distinguished by their chip code sequences (codes) and separate communication channels are created using different codes. Signals from the Node B to the UEs are sent on downlink channels and signals from the UEs to the Node B are sent on uplink channels.\nIn many CDMA systems, a random access channel (RACH) is used for some uplink communications. A RACH is capable of carrying packets of data from multiple UEs. Each packet is distinguishable by a combination of time slot and code. For detection by the Node B, the packets have a sequence which also distinguishes it from other packets. The RACH is a contention-based uplink transport channel which may carry control information from the UE to set up an initial connection with the Node B, for example, to register the UE after power-on to the network or to perform location updates or to initiate a call. Transmissions are sent using repeating frames, each having a plurality of time slots, such as fifteen time slots with only one or two time slots per frame typically dedicated to RACH. When a packet is transmitted over the RACH, it may last for multiple frames. Those frames however, are not necessarily consecutive because a back-off process must be performed between each transmission to control the rate at which UEs access the RACH.\nA UE may attempt a RACH transmission and select a time slot using one of N code identifiers, for example in a TDD CDMA system, one of eight midambles. If no other UE transmits in the same slot with the same midamble and if there is sufficient transmission power, then the UE's RACH transmission succeeds. If another UE transmits in the same slot with the same midamble, then they both fail. This transmission error is known as a collision error. Generally, whenever two or more UEs transmit using the same channel in a wireless system, a collision occurs. Another type of transmission error results when there is insufficient transmission power. The necessary power is generally a function of the channel, the interference, and other PRACH transmissions in the same slot.\nIn some communication systems, such as with a 3GPP system, there is a relatively long delay, on the order of seconds, before which a UE realizes a transmission error has occurred and decides to retransmit the failed packet. The recommended operating condition for the RACH is therefore preferably biased toward having very few collisions or insufficient transmission power errors. The failed packet may be retransmitted on data link layer 2 (L2) or data link layer 3 (L3) depending on the mode of operation.\nThe radio access network has no prior information regarding which RACH codes, or more generally which channels were transmitted. The detection of transmitted transport block sets (TBS) or bursts is performed at the receiver, where the number of UEs that transmitted using the detected code is unknown. In the event of a RACH transmission error, the cause remains unidentified. The error might be the result of a code collision or insufficient transmission power.\nA parameter of dynamic persistence (DP) is defined which is set by the RNC to avoid saturation of the RACH. The DP level (DPL) is broadcast from the Node B to the UEs and the UEs adjust their rate of access to the RACH time slots as a function of DP. A RACH constant value (CV) parameter is defined which is managed at the RNC and is used by the UEs to determine the power of RACH transmissions.\nIn current systems, the DP parameter, RACH CV parameter, and other parameters are set and adjusted in order to avoid collisions and insufficient transmission power errors or, in the alternative, to maintain a predetermined target collision error and target insufficient transmission power error probability. The DP parameter is generated at the Node B and the RACH CV is generated at the RNC.\nA prior art method of controlling these parameters utilizes the number of successful and failed UE transmissions in a time slot for individual system frames. Another prior art method broadcasts these parameters to the UEs, which then adjust their uplink transmission accordingly. It is difficult, however, to appropriately control these parameters because they are separately generated at the Node B and RNC and because the cause of the transmission error remains unknown.\nAccordingly, there exists a need for an improved method of controlling parameters in a contention-based channel wherein the cause of transmission errors is identified and the rate at which such errors occur is identified and controlled by adjusting parameters at the Node B."}
{"text": "Increased interest in the global environmental issues has led to a demand for superior fuel consumption performance in pneumatic tires, along with high wet performance and superior safety performance. As a result, by compounding silica in rubber compositions that form tread portions, heat build-up has been suppressed, rolling resistance reduced, and fuel consumption performance improved and, also, dynamic visco-elasticity characteristics of the tread rubber has been improved which has led to enhancements in wet performance. Efforts have also been made to use silica having a small particulate diameter for the purpose of further increasing the effects associated with the compounding of silica. However, silica has poor affinity with diene rubber and dispersibility tends to be insufficient. Particularly, when the particle diameter of the silica is small, dispersibility worsens and, as a result, the effects of improving low rolling resistance and wet performance by modifying dynamic visco-elasticity characteristics such as the loss tangent (tan δ) of the rubber composition and the like have not been achievable. Additionally, reinforcing effects tend to be lower when compounding silica in a rubber component than when compounding carbon black. As a result, there is a problem in that when dispersibility is poor, wear resistance is insufficient.\nTo resolve this problem, Japanese Unexamined Patent Application Publication No. 2009-091498 as well as patent applications WO/2005/021637 and WO/2003/102053 propose improving the dispersibility of silica by compounding silica in a rubber composition with a terminal-modified solution polymerization styrene butadiene rubber where the terminals are modified by a polyorganosiloxane or the like, thereby reducing heat build-up (tan δ at 60° C.) and enhancing wet grip performance (tan δ at 0° C.). However, the low rolling resistance and wet performance obtained by the technology proposed in Patent Documents 1 to 3 is below that demanded by users and there is a need for further improvement in low rolling resistance and wet performance."}
{"text": "(i) Technical Field\nThe present invention relates to an electronic conference assistance method and an information terminal device employed in an electronic conference system.\n(ii) Related Art\nConventionally, there is available an electronic conference system which has a large-scale readable and writable touch panel display device. Generally, such a touch panel display device is placed so as to be viewed by all participants of the conference and written thereon. Use of the electronic conference system enables a conference of a style, for example, where the participants gather around the touch panel display device, rather than remain seated, to discuss an idea conceived during the conference while writing the idea and so forth on the touch panel display device. The content written on the panel display can be stored intact as a screen image. Also in view of enhancement of conference efficiency, an increasing number of companies are introducing such electronic conference systems.\nMoreover, when such a system is employed, presence of all participants in the conference room where the touch panel display device is installed is not mandatory. That is, when another touch panel display device is installed in a conference room in another location and connected via a network to the electronic conference system in the main location, a remote conference can be realized. This allows a person in a remote location to participate in the conference.\nFurther, when a person who is supposed to participate in the conference but is away from the place where the conference is held as, for example, they are on a business trip connects their own personal computer (PC) to the electronic conference system, they can participate in the electronic conference from any desired place. Still further, when a portable phone is connected to the electronic conference system via a connection line, that person can participate in the conference through audio.\nAs described above, use of an electronic conference system can realize a conference of a style where participants can participate in a variety of manners, not limited to a conventional general conference style in which participants are kept seated and discuss ideas.\nHere, when a remote conference is taking place by connecting the device in the main location to the device used by a conference participant in another location to via a network, basically, data of images captured using a camera or the like in the conference room, in particular data of an image of a person who speaks, is transmitted to other locations. With this arrangement, a participant in the remote location can talk to the person speaking while looking at their image being shown.\nMoreover, a moderator who presides over the conference can check a participant in a conference by looking at the images captured in the respective locations when discussion is carried out between distant locations, and ask the checked participant to present their opinion or encourage them to speak."}
{"text": "Movement of aqueous fluid (e.g., water) within a subterranean formation containing one or more clays (e.g., kaolinite, smectite, illite, chlorite, etc.) often results in reduced permeability (also referred to as “formation damage”) that is adverse to hydrocarbon material (e.g., oil, natural gas, etc.) extraction processes. Such movement can, for example, occur as a result of the introduction of one or more aqueous fluids (e.g., aqueous flooding fluids, other aqueous treatment fluids, etc.) during downhole operations (e.g., drilling operations, acidizing operations, completion operations, flooding operations, hydraulic fracturing operations, squeeze treatment operations, etc.). Particles of various clay minerals (e.g., kaolinite, illite, etc.) may detach from the subterranean formation during the movement of the aqueous fluid and may migrate to and become detained in pore throats of the subterranean formation to impede fluid flow therethrough. Such clay detachment and fluid flow impairment can be exacerbated by higher pH (e.g., lower salinity) aqueous fluids, making switches from lower (e.g., higher salinity) pH fluids to higher pH fluids problematic during downhole operations. Moreover, various other clay minerals (e.g., montmorillonite) may become swelled by the aqueous fluid and hinder fluid flow through adjacent pores. Frequently, movement of aqueous fluid within a subterranean formation containing clay minerals results in both clay-migration-based permeability losses and clay-swelling-based permeability losses.\nVarious conventional clay stabilizers have been used to mitigate subterranean formation permeability losses due to clay migration and/or clay swelling by controlling the charge and electrolytic properties of treatment fluids. Examples of such conventional clay stabilizers include potassium chloride (KCl), sodium chloride (NaCl), zirconium oxychloride (ZrOCl2), sodium hydroxide (NaOH), calcium hydroxide (CaOH), hydroxylated aluminum (Al(OH)3), aluminum salts, zirconium salts, quaternary ammonium salts, and cationic organic polymers. Unfortunately, many such conventional clay stabilizers can be inadequate for prolonged clay stabilization (e.g., being readily removed by acids), can be too large for smaller pores (e.g., contributing to reduced permeability through such pores), can be difficult to handle and/or dispose of, can be environmentally toxic, and/or can be too expensive for practical usage on a commercial scale.\nIt would, therefore, be desirable to have new methods and treatment fluids for stabilizing clay contained within subterranean formations."}
{"text": "Large networks of computers, particularly those which share common hardware components and software applications, are typically administered centrally. As these networks grow in size and complexity, it becomes increasingly difficult for system users (including system administrators) to determine whether and how additions of new software or hardware will affect the existing hardware and software used by the system. Additionally, it becomes more difficult to determine when upgraded versions of existing software arc available and compatible with the existing system. Lastly, the increasing complexity of networked systems makes it more difficult to investigate defects, software dependencies or conflicts, or information critical to the operation of existing software and hardware. In such complex systems, software dependencies are particularly problematic, as many different software applications may scan the same system files for data or software code to enable basic or enhanced functionality or write (and overwrite) data to the same system files. As a result, system users typically treat system files as unalterable xe2x80x9cblack boxes,xe2x80x9d limiting system flexibility.\nExisting system monitoring and management tools, such as the Hewlett-Packard Company\"\"s Openview system or the Computer Associates International, Inc.\"\"s Unicenter systems, employ agents running on target computer systems. These systems rely on an agent or a software process which monitors target computer systems for configuration, diagnostic, frequency of use, and other information which run on the target computer systems to collect configuration and diagnostic information about the target computer systems and relay it to a centralized station. From this station, a system administrator can review the information and take appropriate actions.\nAs described in U.S. Pat. No. 5,933,646 to Hendrickson et al., target computer systems commonly have numerous installed software packages which can be enabled, disabled, installed or removed. A configuration database stores information associated with each of these software packages, including information indicating whether the software package is enabled, the location of the software package within a hierarchical filing structure and a dependency listing (showing the interdependencies between the software package and other software packages on the target computer systems or system files) for each software package. A software manager server supports a software manager graphical user interface (GUI) which permits a system administrator to view the contents of the configuration database. The software manager also carries out administrator-specified changes to the system and amends the configuration database to reflect the changes.\nxe2x80x9cInside TCP/IP,xe2x80x9d by Matthew Flint Arnett et al. (New Riders Publishing, Indianapolis, Ind. 1994) further describes management of networked computer systems. Specifically, it describes the use of simple network management protocol (SNMP) agents to monitor target computer systems. The SNMP agents send error trap messages to a SNMP manager when measured system parameters exceed a threshold value. The SNMP manager may present the error trap signals to a system administrator. The system administrator may set the threshold values at which SNMP traps are triggered.\nThe network of target computer systems may be a local area network (LAN) arranged in a star, bus or ring topology. Network communication priority may be handled according to either the Ethernet or token ring protocols. In an Ethernet LAN, a transmitting station monitors the network transmission channel to determine if it is busy and transmits only after the transmission channel is not busy. In a token ring LAN, communication priority is passed from one station to the next and a station must have priority before it can transmit a message. Network communications can follow the TCP/IP protocol in which messages are divided into discrete xe2x80x9cpacketsxe2x80x9d of information which are addressed to particular receiving addresses.\nThe configurations of target computers in a customer network are typically monitored by a single dedicated system administrator. As changes to the existing software packages are made, old software packages are removed from the target computers or new software packages are added to the target computers, the system administrator must typically change its configuration database in response. Such a process of updating a configuration database incurs significant cost to operating and maintaining the network.\nGenerally, an embodiment of the present invention is directed to a method and system for managing a network of target computers. Another embodiment of the invention is directed to agents for collecting configuration, diagnostic, frequency of use or other information from the target computer system and transmitting the collected information to a central control server. The control server receives the information and accesses relevant information from a database of software information. The control server then formats and transmits this information to the agent. The agent may act on the information directly or may display the information to a user through a management tool GUI.\nAnother embodiment of the present invention is directed to a method of providing system management services to a network of target computers including the steps of enrolling customers to receive computer system management services for a fee, receiving information about the customer network from agents associated with target computer systems, comparing the target computer system information with software and hardware information stored in a database and transmitting that information to the customer."}
{"text": "In an unassisted GPS-type position determination system, subscriber stations determine their own positions from satellite transmissions originating from the GPS-type position determination system, without requesting significant acquisition or calculation assistance from other network entities, for example, dedicated servers. That places significant processing demands on the subscriber stations because of the uncertainty in the timing, position, and frequency of these transmissions, requiring the subscriber stations to expend significant processing resources in searching for and locating these transmissions by, for example, testing large numbers of hypotheses varying the assumed timing, position and frequency of the transmissions. Since the number of hypotheses that must be tested is often staggering, the time required to search for the transmissions can be inordinately long and consume an excessive amount of processing resources, even for subscriber stations with dedicated receiver chains.\nThe uncertainty experienced by the subscriber stations stems from several sources. Assuming GPS positioning, there is first the uncertainty in knowing which of the 32 GPS satellites are visible to the subscriber station. That uncertainty is present because a subscriber station, upon power up or before a position fix is available, has no basis for identifying which signals of these 32 satellites can be usefully received. The useful reception of satellite signals is referred to as an ability of the subscriber station to “see” the satellite emitting the signal, or, in other contexts, as the satellite being “visible” to the subscriber station.\nThis leads to inefficient searching because the subscriber station may waste considerable resources in searching for transmissions from satellites that are not visible to it, and which are therefore not useful for position determination purposes. For example, referring to FIG. 1, while satellites 54a, 54b, and 54c are visible to subscriber station 50 located at position 51 on the earth's surface 52, satellites 56a, 56b, and 56c are invisible to subscriber station 50, as they are located on the other side of the earth. Therefore, it would be wasteful for subscriber station 50 to search for the transmissions from satellites 56a, 56b, and 56c during a position fix attempt.\nIn addition, there is an uncertainty in knowing the timing or phase of the 32 chip PN “gold” codes that are embedded within the individual satellite transmissions. As these codes are circularly shifted versions of one another, the phase of a code uniquely identifies which of the satellites originated the transmission. The phase also reflects the propagation delay caused from transmission from the satellite to the subscriber station. To account for the possible variations in phase, the subscriber station must expend resources in searching over the full range of possible PN codes within a code phase searching window that is large enough to encompass the possible variations.\nMoreover, there is an uncertainty in knowing the relative movement between the subscriber station and the GPS satellites, which typically introduces a Doppler shift of approximately ±4 kHz in the frequency of transmission. To account for the possible variation of frequency introduced by the Doppler shift, the subscriber station must expend resources is searching over the full range of possible transmission frequencies within a frequency searching window that is large enough to encompass the possible variations caused by the Doppler shift.\nFinally, there is the uncertainty in knowing the degree to which the local oscillator (LO) of the subscriber station is out of tune with the GPS carrier frequency. Upon power-up, for example, it is not uncommon for the LO frequency to differ from the GPS carrier frequency by as much as ±5 ppm. Until synchronization between the LO frequency and GPS carrier frequency is achieved, the subscriber station must account for this uncertainty by increasing the size of the frequency search window that is employed.\nEven if the host wireless communications system or GPS-type position determination system eliminating some of this uncertainty by providing timing, positional information, or synchronization to the subscriber station, the processing demands on the subscriber station are often still substantial. For example, a synchronous system, such as a CDMA system, provides the subscriber station with time, and also synchronizes the LO frequency of the subscriber station to the GPS carrier frequency. Although the synchronization substantially reduces the LO frequency uncertainty, for example, from ±5 ppm to ±0.2 ppm, and the timing information allows the subscriber station to determine the position of the GPS satellites (using the GPS almanac or ephemeris data provided by the satellites), the subscriber station is still unable to determine which of the GPS satellites are visible to it, and it is still subject to the frequency uncertainty caused by Doppler shift."}
{"text": "Microprocessors include one or more execution units that perform the actual execution of instructions. Superscalar processors include the ability to issue multiple instructions per clock cycle to the various execution units to improve the throughput, or average instructions per clock cycle, of the processor. However, the instruction fetch and decoding functions at the top of the microprocessor pipeline must provide an instruction stream to the execution units at a sufficient rate in order to utilize the additional execution units and actually improve the throughput. The x86 architecture makes this task more difficult because the instructions of the instruction set are not fixed length; rather, the length of each instruction may vary, as discussed in more detail below. Thus, an x86 microprocessor must include an extensive amount of logic to process the incoming stream of instruction bytes to determine where each instruction starts and ends. Therefore, ways are needed to improve the rate at which an x86 microprocessor can parse a stream of indistinct instruction bytes into distinct instructions."}
{"text": "1. Field of the Invention\nThe present invention relates to a switching power supply apparatus in which at least two output voltages are produced by means of one transformer.\n2. Description of the Related Art\nGenerally, in a switching power supply apparatus having the configuration in which two output voltages are obtained by means of one transformer, only one output voltage is detected to control a first switching element connected to the primary of the transformer. FIG. 1 is a circuit diagram of a conventional switching power supply apparatus configured as described above.\nA transformer T is provided with a primary winding T1, secondary windings T2 and T3, and a drive winding T4. A switching element Q1 (hereinafter, referred to as a first switching element) is connected in series with the primary winding T1. A control circuit CT controls the on-time of the first switching element Q1. Output from the secondary winding T2 is rectified by a rectification diode Ds1 and smoothed by a capacitor to be output as a first output Vo1. Moreover, output from the other secondary winding T3 is rectified by a rectification diode Ds2 and smoothed by a capacitor to be output as an output Vo2. In this example, the control circuit containing the drive winding T4 of the transformer T controls the on-time of the first switching element Q1, and moreover, causes the first switching element Q1 to oscillate autonomously. With this configuration employed, when the first switching element Q1 is on, an input voltage Vin is applied to the primary winding T1, making input current flow so that energy is stored in the transformer T. Furthermore, when the first switching element Q1 is off, the energy stored in the transformer T is released as output current via the secondary windings T2 and T3. Thus, the device is configured as an energy storage type switching power supply apparatus .\nIn the control circuit CT, output voltage from the drive winding T4 is delayed to some degree by a resistor R4 and a capacitor C3 to be applied to the control terminal of the first switching element Q1, and moreover, is applied to a time constant circuit comprising a resistor R1 and a capacitor C1, so that a transistor Tr1 is turned on after a constant time-period, causing the first switching element Q1 to turn off. When the first switching element Q1 turns off, energy stored in the transformer T is released as electric current. When the release of the energy is completed, reverse voltages are applied to the rectification diodes Ds1 and Ds2, respectively. The capacitive impedance, equivalent from the standpoint of these rectification diodes, and the winding inductor of the transformer T resonate, and with voltage generated in the drive winding T4 of the first switching element Q1, voltage is applied to the control terminal of the first switching element Q1, so that the first switching element Q1 turns on again.\nMoreover, an output voltage detection circuit DT is provided on the output side of the rectification circuit comprising the rectification diodes Ds1 and Ds2, and capacitors. The output voltage detection circuit DT detects only output voltage with respect to the first output Vo1. That is, the voltage of the first output Vo1 is detected by resistors R2 and R3. The voltage divided by the resistors R2 and R3 is input to a voltage comparison terminal Vr as a comparison voltage. A series circuit comprising a photodiode PD1, a shunt regulator ZD1, and a resistor 5 is connected between Vo1 and GND. The above detection voltage is input to the voltage comparison terminal (reference terminal) Vr of the shunt regulator ZD1. A phototransistor PTr1 arranged in opposition to the photodiode PD1 is connected between the base and collector of the transistor Tr1 in the above-described control circuit CT.\nWith the above arrangement, when the first output Vo1 is increased, the input voltage to the shunt regulator ZD1 is increased. Then, since the inflow current to the photodiode PD1 is increased, the transistor Tr1 is turned on earlier via the operation of the phototransistor PTr1 in the control circuit CT. As a result, the on-time of the first switching element Q1 becomes shortened, and thereby, the first output voltage Vo1 is reduced. In this way, the output voltage of the first output Vo1 is monitored by the output voltage detection circuit DT, and a control signal corresponding to the detection voltage is formed and fed back to the control circuit CT, whereby the first output voltage Vo1 can be stabilized. Moreover, since the first output Vo1 is stabilized, the second output voltage Vo2 is stabilized to some degree.\nHowever, the above-described conventional switching power supply apparatus shown in FIG. 1 has the problem that, though the output voltage of the first output voltage Vo1 to which the output voltage detection circuit is directly connected can be stabilized with high accuracy, a sufficiently high voltage accuracy can not be obtained for the output voltage Vo2 other than the output voltage Vo1.\nTo solve this problem, in some cases, a voltage stabilization circuit such as a series regulator circuit or the like is inserted in the second output voltage Vo2 circuit, or a dummy resistor is used. With such circuits, problems are caused such as an increase in number of parts, reduction of circuit efficiency, temperature rise of the power supply apparatus, and so forth.\nAccordingly, it is an object of the present invention to provide a switching power supply apparatus in which, when at least two output voltages are provided, each of the output voltage accuracies can be stabilized at a predetermined control ratio by detecting the respective output voltages.\nThe switching power supply apparatus of the invention comprises a transformer having a primary winding and at least two secondary windings, a first switching element connected in series with the primary winding, a control circuit for controlling the output from the first switching element by control of the on-time thereof, a rectification circuit for rectifying at least two outputs from the secondary windings, and an output detection circuit for detecting the output voltages, and feeding back the output voltages as a control signal for the on-time to the control circuit, wherein the output voltage detection circuit comprises a control signal formation section in which the control signal for the on-time is formed, corresponding to voltage at a voltage comparison terminal, and plural voltage detectors connected between the at least two outputs of the secondary windings and the voltage comparison terminal, respectively.\nIn this switching power supply apparatus, output voltages from the respective outputs, generated by at least two secondary windings, are detected by the voltage detectors connected to the outputs, respectively, and are input together to the voltage comparison terminal. The extents of the influences of the variations in the respective outputs onto the voltage comparison terminal will be considered below. The extents of the influences, if resistors are used for the detection, are varied, depending on the resistances. Accordingly, the plural output voltages can be controlled at an optional ratio by designing the resistances of the voltage detection resistors corresponding to their specifications.\nIn this way, the plural voltage detectors are connected between at least two outputs of the secondary windings and the voltage comparison terminal, and the control signal for the on-time of the first switching element connected to the primary winding is formed, corresponding to voltage at the voltage comparison terminal. Therefore, the respective output voltages can be stabilized at a desired control ratio.\nAccording to an aspect, for the plural voltage detectors, at least one Zener diode is used.\nBy appropriate selection of the Zener diode, the respective output voltages from the circuit can be stabilized at the above-described desired control ratio only when the output voltage from the circuit connected to the Zener diode exceeds a predetermined voltage. Thus, the output voltages can be suppressed from increasing. Moreover, when the output voltage from the circuit connected to the Zener diode is less than the predetermined voltage, only the output voltage from the circuit not connected to the Zener diode is controlled for stabilization.\nAccording to another aspect, the secondary winding comprises at least two secondary windings, the rectification output terminal of a predetermined secondary winding is connected to one end of the other secondary winding, and current through the other end of the other secondary winding is rectified, and is output.\nIn this embodiment, no influences are exerted due to of variations in voltage, caused by variations in current flowing in the rectification diode for the predetermined secondary winding. For this reason, correspondingly, the voltage accuracy of the other secondary winding is improved.\nAccording to still another aspect, when the first switching element is on, input voltage is applied to the primary winding, causing current to flow so that energy is stored in the transformer, and when the first switching element is turned off, the energy stored in the transformer is released from the secondary windings.\nIn this embodiment, the switching power supply apparatus comprises a flyback type. Thus, it is not necessary to provide a choke coil or the like on the secondary side. Accordingly, a switching power supply apparatus having a small size, a high accuracy, and a high stability can be provided.\nAccording to still another aspect, the switching power supply apparatus further comprises an inductor connected in series with the primary winding, and a series circuit comprising a capacitor and a second switching element, connected in parallel to the series circuit comprising the inductor and the primary winding, wherein the control circuit turns the first and second switching elements on and off, alternately, so as to sandwich a time-period when both of the switching elements are off, and controls the on-time of the switching elements, whereby the outputs therefrom are controlled.\nSuch a switching power supply apparatus in which the primary of the transformer is configured as described above is disclosed in U.S. Pat. No. 6,061,252 and Japanese Unexamined Patent Publication No. 11-187664, both of which are assigned to the assignee of the present invention, and the disclosure of which are hereby incorporated by reference.\nIn this switching power supply apparatus, when the first switching element turns off, energy stored in the inductor connected in series with the primary winding is released as charging current into the capacitor. Then, directly after this, the second switching element turns on, and discharging is carried out based on the charge potential of the capacitor. With the discharging current, energy is stored in the primary winding of the transformer and the inductor. When the second switching element turns off after a predetermined time-period, the energy stored again in the inductor L flows via the primary winding and the input power source. In this operation, the inductor connected in series with the primary winding includes the leakage inductance of the transformer. Accordingly, generation of surge, caused by the leakage inductance at switching, can be prevented. Moreover, the discharging current, generated when the second switching element Q2 is on, becomes an resonant current. This is reflected by the secondary. That is, the secondary winding current output takes a part of the sinusoidal waveform starting from a zero voltage (mountainous waveform), so that surge in a leading edge can be practically neglected.\nSince the current surge or the like is suppressed as described above, the output voltages can be prevented from increasing due to the surge current or the like. As a result, the voltage accuracy of an output produced at a small control ratio or a non-controlled output can be improved. Especially, when said load having a large control ratio is heavy, and the load having a small control ratio is light, conventionally surge current causes the output having a small control ratio to rise in voltage. However, this is considerably improved by this arrangement.\nAccording to still another aspect, the switching power supply apparatus further comprises an inductor connected in series with the primary winding, and a series circuit comprising a capacitor and a second switching element, connected in parallel to the first switching element, wherein the control circuit turns the first and second switching elements on and off, alternately, so as to sandwich a time-period when both of the switching elements are off, and controls the on-time of the switching elements, whereby the outputs therefrom are controlled.\nWith this structure, the switching power supply apparatus carries out the same operation as expalained above.\nAccording to still another aspect, the leakage inductance of the transformer is used.\nIn this embodiment, the inductor consists of a leakage inductance itself. This can reduce the number of parts, since it is not necessary to provide the inductor as a separate part.\nAccording to still another aspect, the control circuit comprises a drive winding provided in the transformer to drive the first and second switching elements, respectively, and a control section provided with a time constant circuit for providing on-off signals to the control terminals of the first and second switching elements at a predetermined timing by use of a voltage substantially proportional to the voltage of the primary winding and generated in the drive winding, whereby the first and second switching elements oscillate autonomously.\nIn this embodiment, the first switching element and the second switching element are autonomously operated. Thus, oscillation IC\"\"s or the like are not needed. The number of parts can be significantly reduced. Moreover, the first and second switching elements can be easily turned on and off alternately so as to sandwich a time-period when both of them are off, and the on-time of these switching elements can be simply controlled. The loss and breaking of elements, caused by the short-circuit current which flows when the two switching elements are simultaneously turned on, can be prevented.\nAccording to still another aspect, the control circuit includes a rectification diode, and a capacitive impedance connected in parallel to the rectification diode.\nSince the first switching element in the primary operates as a switch, surge-voltage and surge-current are produced in the output of the secondary. The capacitive impedance is connected in parallel to the rectification diode, which enables the voltage surge to be absorbed. Moreover, charges are supplied to the output via the capacitive impedance, so that the affects of the voltage-drop of the rectification diode can be reduced.\nAccording to still another aspect, the rectification circuit includes an inductive impedance connected in series with the rectification diode.\nIn this embodiment, the inductive impedance is connected in series with the rectification diode, Thus, especially, current surge can be prevented.\nFor the purpose of illustrating the invention, there is shown in the drawings several forms which are presently preferred, it being understood, however, that the invention is not limited to the precise arrangements and instrumentalities shown.\nOther features and advantages of the present invention will become apparent from the following description of the invention which refers to the accompanying drawings."}
{"text": "In general, it is preferable to use a base in a heat-developable light-sensitive material for the purpose of accelerating the heat development of said material, and it is necessary to incorporate said base in said light-sensitive material in the form of a base precursor for increasing the stability of said material. In order to use such base precursor practically, which means herein a compound capable of being decomposed under heat to be able to release a basic component, it is necessary that said base precursor must have both stability at normal temperature and rapid decomposability under heat.\nVarious kinds of conventional base precursors have heretofore been known including, for example, ureas as described in U.S. Pat. No. 2,732,299 and Belgian Patent No. 625,554; urea or ammonium salts of urea and weak acid as described in Japanese Patent Publication No. 1699/65; hexamethylene-tetramines or semicarbazides as described in U.S. Pat. No. 3,157,503; combinations of triazine compound and carboxylic acid as described in U.S. Pat. No. 3,493,374; dicyan-diamide derivatives as described in U.S. Pat. No. 3,271,155; N-sulfonylureas as described in U.S. Pat. No. 3,420,665; amine-imides as described in \"Research Disclosure\" (1977), RD No. 15776; and salts of decomposable acids such as trichloroacetic acid as described in British Patent No. 998,949.\nHowever, image forming materials containing such conventional base precursors have serious defects. In particular, such conventional base precursors do not fully satisfy the properties desired both with respect to high stability during preservation at normal temperature (e.g., 0.degree.-30.degree. C.) and rapid decomposability during development treatment under heat. Therefore, a high and sufficient image density cannot be obtained, or the S/N (signal/noise) value of the formed image extremely decreases, as the base component tends to be released from said base precursor during preservation thereof, and these phenomena are extremely serious problems."}
{"text": "The present invention relates to an X-ray image intensifier.\nAs X-ray image intensifiers (to be referred to as \"I.I.\"s hereinafter), a general-purpose single visual field type I.I. and a high-grade variable visual field type I.I. are frequently used. In general, an I.I. comprises a vacuum housing which includes a substantially cylindrical outer casing, and an X-ray entrance window and an X-ray exit window which are arranged to close two ends of the outer casing. In the vacuum housing, input and output surfaces are arranged along the entrance and exit windows, respectively, and a focusing electrode constituting an electronic lens is located between the input and output surfaces. The I.I.s are classified into the single visual field type and variable visual field type due to differences in the number and arrangement of focusing electrodes, and the like. In the case of a variable visual field type I.I., when a voltage distribution to the focusing electrodes is switched, an output visual field image can be enlarged like, a normal visual field, a second visual field, a third visual field,....\nThe input surface has a base and a phosphor screen formed on the base, and has an arcuated circular shape.\nIn U.S. Pat. No. 3,716,713, the thickness of the phosphor screen is increased from its center toward the periphery, and is maximized at the periphery.\nAccording to an I.I. disclosed in Japanese Patent Disclosure No. 53-102663, the phosphor screen has the same arrangement as that in the above U.S. Pat. No., and the base has a mosaic structure having a large number of grooves for effecting a light guide function.\nAccording to an I.I. disclosed in Japanese Patent Disclosure No. 59-207551, the thickness of the phosphor screen is decreased from its center toward the periphery, and X-ray optical path lengths passing through the phosphor screen are adjusted to be equal to each other at the center and the periphery of the phosphor screen.\nIn the I.I.s having the above-mentioned arrangements of the input surfaces, the characteristic of an image obtained at the output surface, in particular, a luminance distribution characteristic, is such that a luminance is high at the center of the image and is gradually decreased toward the periphery. Therefore, a luminance distribution curve obtained as a result of measurement along the diameter of an image becomes a quadratic curve. In the variable visual field I.I., the same luminance distribution characteristic is obtained either in a normal visual field operation or in an enlarged visual field operation.\nThe reason for the above-mentioned luminance distribution can be considered as follows.\nIn the I.I.s disclosed in U.S. Pat. No. 3,716,713 and Japanese Patent Disclosure No. 53-102663, in order to prolong an X-ray passage distance in the phosphor screen, which influences light emission, so as to compensate for a quantity of light emitted from the phosphor screen, the thickness of the peripheral portion of the phosphor screen is increased. However, a portion between the intermediate portion and periphery of the phosphor screen cannot provide a similar effect upon increase in thickness, and, to the contrary, the luminance of the periphery of an image is decreased. This is because an excessive increase in thickness at the peripheral portion of the phosphor screen does not contribute to light emission of the phosphor by means of X-rays but degrades a transmittance of X-rays.\nIn Japanese Patent Disclosure No. 59-207551, in order to obtain a constant passage distance of X-rays at respective positions in the phosphor screen, the thickness of the phosphor screen is decreased at a given rate from its center toward the periphery. However, in order to obtain a theoretical luminance, the phosphor screen must be formed to have a uniform structure and a uniform emission intensity distribution. If these conditions cannot be satisfied, the luminance at the peripheral portion of an image, in particular, an area shifted from the center of the image toward the periphery by a distance 80 to 95% of an effective image diameter, is considerably decreased as compared to the above two prior arts.\nWhen the I.I.s having the above luminance distribution characteristic are used, the following problems are posed. In the distribution characteristic, the luminance at the center of an image is high and is decreased toward the periphery. When the I.I. is coupled to an optical system, a luminance difference between the center and the periphery of the image is emphasized due to an operation of the optical system. For this reason, a dark portion at the peripheral portion of the image has degraded discriminating ability of an object, and cannot be used for observing an object. Therefore, a virtual image area is decreased. When an object is observed upon clinical examination, a contour image of the object must be confirmed. However, when the effective image area is small as described above, the I.I. must be moved stepwise so that a portion to be observed is located at the center of the image. For this reason, the observation requires a long time, and an X-ray irradiation time is also prolonged. For example, when an observation is performed using a TV fluoroscopic imaging method, the entire object, i.e., the entire image, must be scanned, and this requires still more time.\nIn the enlarged visual field operation mode, e.g., in the second visual field operation mode, the luminance distribution characteristic of an output image is such that the center of an image is bright and the peripheral portion thereof is dark as in the normal visual field operation mode. In any visual field operation mode, an area of an input visual field is changed, but an image area which can be observed is almost not changed. For this reason, when the enlarged visual field operation is performed in order to microscopically observe the object after the contour image of the object is confirmed, the I.I. must be moved to locate the object at the center of image. If the object is a moving body, and is moved to the peripheral portion of an output image, the object cannot be discriminated since the luminance of the peripheral portion is low.\nSince the luminance distribution characteristic is not changed in the enlarged visual field operation mode, a low luminance portion is moved upon switching of visual fields. The object is often out of sight upon switching of the visual fields, and the I.I. must be moved to confirm the object at that time. For example, upon clinical examination wherein a change in object must be immediately judged, such as blood vessel imaging, the lack of necessary data and the complicated operations as described above may cause serious problems."}
{"text": "This invention relates to an ignition system and more particularly to an ignition system for an internal combustion engine. The invention also relates to an alternative spark-plug, a drive circuit for a spark-plug and associated methods.\nIt is known that an ignition system for a vehicle comprises a plurality of distributed spark-plugs connected by respective high voltage power cables to a remote and central high voltage generation means. In a known capacitor discharge ignition system, the high voltage generation means comprises a capacitor connected with a power switching device, such as an SCR switch, in series with a primary winding of a transformer. A secondary winding is connected to the high voltage cables. In use, when a piston of the engine reaches a predetermined position, the power switching device is switched to the closed state. Energy in the capacitor is then transferred to the primary winding resulting in a much higher voltage on the secondary, because of the secondary to primary winding ratio. Once the voltage on the secondary reaches the breakdown voltage of a spark-gap between spark electrodes of the plug, a plasma discharge is created between the spark electrodes.\nIn the known systems, the switching circuit restricts the minimum inductance of the transformer that can be used. The restricting factors are the maximum current rating of the switch, Im, the switching speed of the switch ts, the switching voltage of the switch, Vs, and the cost of the switch. These limitations result in a very high secondary winding inductance, which has several drawbacks including cost. The large inductance normally requires kilometres (ten thousands of windings) of thin copper wire, which is expensive. The systems are inefficient in that the kilometres of thin copper wire have a resistance of a few kilo-ohms. To transfer enough energy for a reliable spark, a large amount of extra energy is required for each spark. Due to the large amount of energy that must be handled as well as the large amount of copper needed, the systems are bulky. The energy loss due to the copper resistance, heats the transformer. This places a severe limit on the maximum amount of energy that can be transferred to the spark and also affects the placement of the transformer for cooling. The fuel efficiency, completeness of combustion, combustion time, exhaust cleanliness and variability in cycle-to-cycle combustion are limited. Because the transformer is large and heats up, it is normally positioned a distance away from the engine. This requires high voltage cables between spark-plugs and the transformer. These high voltage cables generate a large amount of electromagnetic radiation, which may influence other electronic equipment. In order to eliminate the high voltage cables, coil-on-plug systems which comprise an ignition coil at each spark-plug are used. Because these coils are very close to the engine, normally with very little air flow around them, they overheat easily, which makes them unreliable.\nSome ignition coils having a very low secondary resistance have been suggested. This is accomplished by using a magnetic path having a high permeability, to reduce the number of windings while keeping the inductance high enough for the switching circuit. The disadvantage of this approach is that the high permeability magnetic material saturates easily and that a large core is therefore required.\nSome other ignition systems have a second energy transfer path on the secondary side. They all have the disadvantage that the energy must either go through the secondary winding or through a semiconductor device. If the energy goes through the secondary winding, the transfer is very inefficient due to the high winding resistance. On the other hand, the semiconductor device must be a high voltage (normally above 30 kV), high current (normally above 1 A) device. These devices are expensive and also result in energy loss.\nAnother disadvantage of all these systems is that the self-resonance frequency of the secondary winding is low (typically less than 20 kHz). The low self-resonance frequency is due to the long length of secondary wire and the large secondary winding inductance. When the secondary winding is connected in a secondary side circuit, the resonance frequency of the secondary side circuit is even lower than the self-resonance frequency of the secondary winding, due to the spark-plug and cable capacitance. Because of the low secondary resonance frequency, it takes some tens of microseconds to charge the spark-plug or electrode capacitance to a breakdown voltage and also some tens of microseconds to dissipate the remaining secondary energy. This limits the number of successive pulses that can be generated in multiple spark ignition systems, which limits the amount of energy that can be delivered during ignition. The efficiency and amount of energy transferred in some ignition systems are increased by placing a capacitor in parallel with the spark-plug. In these systems the secondary resonance frequency will be even lower. Even in systems where an optimal spark time is calculated (as discussed below), the spark cannot be controlled to within a few tens of microseconds. At 6000 rpm, this inaccuracy is larger than one degree in engine rotation.\nIt is a known technique to use the spark-plug to measure the current in or resistance of the ionized gas after ignition to gain information about the gas temperature, pressure or composition after combustion. This information is then used as one of the inputs to an engine management system to calculate an average optimal spark time. Because of the high loss of the ignition transformer, the measurement must be done on the secondary side of the transformer, which makes the secondary side circuit complex.\nDue to cycle-to-cycle variations, the average optimal spark time can be quite different from the optimal spark time for a single cycle. Although there are a number of techniques available to measure the conditions inside the combustion chamber before ignition, none of them are widely used because they all require extra access points to the combustion chamber, are expensive, most have low reliability and are complex.\nWhen using the spark-plug for measurements, the low secondary resonance frequency therefore limits the measuring frequency after ignition and also makes it very difficult, if not impossible, to measure gas properties before ignition."}
{"text": "Depending on the applications envisaged, thermal stabilization during shaping or during the use of the corresponding products is carried out in various ways.\nFrench Patent FR-A-2,297,227 describes PVC compositions stabilized effectively by metal-organic salts such as, salts of zinc, calcium and barium and .beta.-diketones.\nThese compositions also contain common additives such as epoxidized oils, lubricants, plasticizers or impact strength enhancing agents.\nFor some applications, especially those for which the chlorinated polymer has to be plasticized to only a small extent, if at all, the customary presence of relatively large amounts of epoxidized oil, such as epoxidized soybean oil or epoxidized linseed oil, in the polymer, tends to lower the softening point of the polymer.\nIn fact, the elimination or a significant decrease in the content of these epoxidized oils decreases the thermal stability of the polymer compositions.\nThis manifests itself, during thermo-forming of the compositions, in a substantial yellowing of the polymer, which is unacceptable for applications requiring transparency and a colorless or only slightly colored appearance of the shaped article, for example, in the case of PVC containers."}
{"text": "The present invention relates generally to microlasers and associated fabrication methods and, more particularly, to Q-switched microlasers and associated fabrication methods.\nModern electro-optical applications are demanding relatively inexpensive, miniaturized lasers capable of producing a series of well-defined output pulses. As such, a variety of microlasers have been developed which include a microresonator and a pair of at least partially reflective mirrors disposed at opposite ends of the microresonator to define a resonant cavity therebetween. The microresonator of one advantageous microlaser includes an active gain medium and a saturable absorber that serves as a Q-switch. See, for example, U.S. Pat. No. 5,394,413 to John J. Zayhowski, which issued on Feb. 28, 1995, the contents of which are incorporated in their entirety herein. By appropriately pumping the active gain medium, such as with a laser diode, the microresonator will emit a series of pulses having a predetermined wavelength, pulse width and pulse energy.\nAs known to those skilled in the art, the wavelength of the signals emitted by a microlaser is dependent upon the materials from which the active gain medium and the saturable absorber are formed. In contrast, the pulse width of the laser pulses emitted by a conventional microlaser is proportional to the length of the resonator cavity. As such, longer resonator cavities will generally emit output pulses having greater pulse widths. Further, both the pulse energy and average power provided by a microlaser are proportional to the pulse width of the pulses output by the microlaser. All other factors being equal, the longer the microresonator cavity, the longer the pulse width and the greater the pulse energy and average power of the resulting laser pulses.\nConventional microlasers, such as those described by U.S. Pat. No. 5,394,413, are end pumped in a direction parallel to the longitudinal axis defined by the resonator cavity. In this regard, the longitudinal axis of the microresonator cavity extends lengthwise through the resonator cavity. Since the resonation cavity is generally a rectangular solid, the longitudinal axis is oriented so as to be orthogonal to the pair of at least partially reflective mirrors that define the opposed ends of the resonant cavity. As such, conventional microlasers are configured such that the pump source provides pump signals in a direction perpendicular to the at least partially reflective mirrors that define the opposed ends of the resonant cavity. The effective length of the resonator cavity is therefore equal to the physical length of the resonator cavity.\nWhile the microlaser can be fabricated such that the resonator cavity has different lengths, a number of factors contribute to generally limit the permissible length of the resonator cavity. See, for example, U.S. Pat. No. 5,394,413 that states that the resonator cavity, including both the saturable absorber and the gain medium, is preferably less than two millimeters in length. In particular, a number of electro-optical applications require microlasers that are extremely small. As such, increases in the length of the resonator cavity are strongly discouraged in these applications since any such increases in the length of the microresonator cavity would correspondingly increase the overall size of the microlaser.\nIn addition, the length of passively Q-switched microlasers is effectively limited by the requirement that the inversion density must exceed a predetermined threshold before lasing commences. As the physical length of the resonator cavity increases, greater amounts of pump energy are required in order to create the necessary inversion density for lasing. In addition to disadvantageously consuming more power to pump the microlaser, the increased pumping requirements create a number of other problems, such as the creation of substantially more heat within the microlaser which must be properly disposed of in order to permit continued operation of the microlaser. In certain instances, the heat generated within the microlaser may even exceed the thermal capacity of the heat sink or other heat removal device, thereby potentially causing a catastrophic failure of the microlaser.\nSince the pulse width and correspondingly the pulse energy and average power of the pulses output by a microlaser cavity are proportional to the length of the resonator cavity, the foregoing examples of practical limitations on the length of the resonator cavity also disadvantageously limit the pulse width and the corresponding pulse energy and average power of the pulses output by conventional microlasers. However, some modem electro-optical applications are beginning to require microlasers that emit pulses having greater pulse widths, such as pulse widths of greater than 1 nanosecond and, in some instances, up to 10 nanoseconds, as well as pulses that have greater pulse energy, such as between about 10 xcexcJ and about 100 xcexcJ, and greater average power, such as between 0.1 watts and 1 watt. As a result of the foregoing limitations on the length of the resonator cavity and the corresponding limitations on the pulse widths, pulse energy and average power of the pulses output by the conventional microlasers, conventional microlasers do not appear capable of meeting these increased demands.\nA microlaser is therefore provided according to one embodiment of the present invention that is capable of supporting a zig-zag resonation pattern in response to pumping of the active gain medium so as to effectively lengthen the microresonator cavity without having to physically lengthen the microresonator cavity. As such, the microlaser of these embodiments can generate pulses having greater pulse widths and correspondingly greater pulse energies and average power levels than the pulses provided by conventional microlasers of a similar size.\nAccording to the present invention, the microlaser includes a microresonator having an active gain medium and a Q-switch, such as a passive Q-switch proximate to and, in one embodiment, immediately adjacent to the active gain medium. In advantageous embodiments, the active gain medium and the Q-switch are integral such that the microresonator may be a monolithic structure. The microresonator extends lengthwise between opposed end faces. The microlaser also includes first and second reflective surfaces disposed proximate respective ones of the opposed end faces to define a microresonator cavity therebetween. While the first and second reflective surfaces can be coated upon respective ones of the opposed end faces of the microresonators, the first and second reflective surfaces can also be formed by mirrors that are spaced from respective ones of the opposed end faces. The microlaser can also include a pump source for introducing pump signals into the active gain medium via at least one of the end surfaces of the microresonator such that the zig-zag resonation pattern is established within the microresonator cavity.\nIn one advantageous embodiment, the opposed end faces are each disposed at a nonorthogonal angle xcex1, such as between about 30xc2x0 and about 45xc2x0, relative to a line perpendicular to a longitudinal axis defined by the microresonator cavity and extending between the opposed end faces. In one embodiment, the opposed end faces are each disposed at the same nonorthogonal angle xcex1 relative to the longitudinal axis such that the opposed end faces are parallel. In another embodiment, the opposed end faces are oriented in opposite directions by the same nonorthogonal angle xcex1. As a result of the nonorthogonal relationship of the opposed end faces, the microlaser of either embodiment is capable of supporting the zig-zag resonation pattern in response to pumping of the active gain medium via at least one of the end surfaces of the microresonator.\nBy supporting the zig-zag resonation pattern, the effective length of the microresonator cavity is increased relative to conventional microlasers having substantially the same physical size that do not support a zig-zag resonation path. In this regard, the effective length of the microresonator cavity of the present invention is the length of the zig-zag resonation path established by the microlaser which is significantly longer than the linear resonation paths established by conventional microlasers that extend parallel to the longitudinal axis of the resonator cavity. As such, the microlaser of the present invention can emit pulses having a longer pulse width and correspondingly greater pulse energies and average power levels than the pulses emitted by conventional microlasers of the same physical size.\nIn order to permit the pump signals to be received by the active gain medium without being reflected from the end face, the microlaser can include an antireflection coating on the end face through which the pump signals are delivered for permitting pump signals having a predetermined range of wavelengths to be received by the active gain medium. The microresonator also generally includes a plurality of side surfaces extending between the opposed end faces. In order to further facilitate resonation within the microresonator cavity, the plurality of side surfaces can be roughened, such as by grinding, to thereby diffuse light.\nIn order to permit the microlaser to emit signals of a predetermined lasing wavelength via one of the opposed end faces, the first reflective surface is preferably highly reflective for laser signals having the predetermined lasing wavelength. In contrast, the second reflective surface is preferably only partially reflective for laser signals having the predetermined lasing wavelength. As such, the microlaser can emit laser pulses having the predetermined lasing wavelength via the second reflective surface.\nIn one embodiment, the microlaser also includes a heat sink upon which at least the microresonator is mounted and a housing in which at least the microresonator is disposed. In this embodiment, the housing includes a window through which laser signals generated by the microresonator are emitted."}
{"text": "1. Field of the Invention\nThis invention relates generally to light emission devices incorporating multiple solid state emitters (e.g., light emitting diodes (LEDs)) and/or multiple phosphors as emissive components.\n2. Description of the Related Art\nIn the illumination art, a variety of approaches have been employed to produce light of desired spectral character.\nLEDs have come into widespread usage as a result of their advantages, which include small size, long life, low energy consumption, and low heat generation.\nU.S. Pat. No. 6,513,949 issued Feb. 4, 2003 describes a hybrid lighting system for producing white light, including at least one LED and a phosphor-LED, in which the color and number of the LEDs and/or the phosphor of the phosphor-LED may be varied.\nU.S. Pat. No. 6,600,175 issued Jul. 29, 2004 describes a light emitting assembly including an LED emitting a first, relatively shorter wavelength radiation, and a down-converting luminophoric medium, e.g., phosphoric medium, that in exposure to such first radiation, responsively emits a second, relatively longer wavelength radiation.\nWhite LED devices have been commercialized that utilize a blue LED and a YAG phosphor (Y3Al5O12 doped with cerium) that partially absorbs blue radiation (centered on 470-480 nm) from the blue LED and emits light having a broad wavelength range with a dominant yellow characteristic (centered on ˜550-600 nm).\nThe commercially available LED/phosphor devices for production of white light do not provide high conversion efficiency color rendering in various spectral regimes of interest. For example, in many applications, consumers prefer white light having color (as quantifiable by color temperature and color rendering index values) that matches sunlight, conventional incandescent bulb light, or fire light such as candle light.\nThere is accordingly a continuing need in the art for efficient LED/phosphor illumination systems producing light having a color rendering that closely matches a predetermined spectral distribution."}
{"text": "1. Field of the Invention\nThis invention relates to an improvement in a collapsible and readily mailable accessory for desk telephones which is easily attachable to serve as a memorandum pad support and writing instrument holder and which also includes a support for a calender in upstanding position and may have a holder for a personal directory in either open or closed position.\n2. Description of the Prior Art\nThis invention is an improvement of my quickly attachable note pad accessories for desk telephones disclosed in my prior U.S. Pat. No. 4,004,112 and corrects an instability and possible displacement of the note pad supporting platform from its intended position which is occasionally encountered when the telephone hand set is removed from the cradle or when the platform is inadvertently struck by the user. The platform of my prior accessory is constructed for manufacture by stamping from sheet metal or by molding of plastic which is then assembled with the wire brace. This results in a costly procedure for a product which is intended for merchandising primarily through the advertising and premium trade."}
{"text": "It has been found that various chemical compounds, such as, for example, nitric oxide (NO), administered during a patient inspiratory phase may provide beneficial effects.\nFor example, NO presents some lung vasodilator properties that may be helpful for respiratory distress conditions such as respiratory distress syndrome of newborn.\nApparatus for delivering such gaseous chemical compounds have therefore been designed to deliver the compounds during the patient's inspiratory phase.\nOne such apparatus is described in Canadian Patent Application No. 2,106,696, filed on Sep. 22, 1993 and published on Mar. 25, 1994 and naming Robert Briand and Marie-Hélène Renaudin as inventors. In this document, Briand et al. describe an apparatus for delivering controlled doses of NO to the respiratory system of the patient without conventional pre-mixing of the NO with oxygen supplied by a ventilator device. The apparatus therefore includes means for detecting the beginning of a patient inspiratory phase and to open an electromagnetic valve for a predetermined duration to supply a controlled dose of NO. The duration and the pressure of the NO supplied dose is adjusted so as to obtain the desired NO concentration with respect to the average inhalation volume of the patient. The NO dose supplied is therefore not directly related to the inhalation volume of the patient. Of course, there is no NO injection during the expiration phase.\nA major drawback of the apparatus described by Briand et al. is the automatic opening of the electromagnetic valve for a predetermined duration each time the beginning of an inhalation phase is detected. Indeed, if the patient repetitively draws short breaths, harm may be caused by the high concentration of NO injected to the patient.\nIn an article entitled: “Comparison of two administration techniques on inhaled nitric oxide on nitrogen dioxide production”, published in Canadian journal of Anaesthesiology 1995, Vol. 42: 10, pages 922–927, Dubé et al. describe an injection system for delivering NO during inspiratory phase. In this injection system, an electronic circuit detects the beginning and the end of each inspiration by processing a flow signal supplied by a ventilator. At the beginning of the inspiratory phase, the electronic circuit opens a solenoid valve and NO is injected into the respiratory line. At the end of the inspiratory phase, the electronic circuit closes the solenoid valve and the injection of NO is stopped.\nFIG. 1 of the appended drawings is a graph of the inspiratory gas flow 20 vs time for a conventional ventilator when the ventilator is in a first mode. When it is in this mode, the flow of inspiratory gas is constantly delivered for a predetermined duration (inspiratory phase 22) and the patient then expires (expiratory phase 24). In the injection system proposed by Dubé et al., when the gas flow reaches a predetermined threshold level 26, a solenoid valve is open, delivering NO to the patient. The line 28 illustrates the injected flow of NO in the inspiration circuit over time. It is to be noted that the scale is different for the flow of inspiratory gas 20 and the flow 28 of NO. Indeed, line 28 illustrating the flow of NO is shown scaled up for illustrative purposes.\nSince the solenoid valve used by Dubé et al. is of the type fully open/fully closed, the flow 28 of NO is constant when the valve is open. As can be seen from FIG. 2, the concentration 29 of NO is essentially constant over time during the inspiratory phases. When the inspiratory gas flow 20 falls below the threshold level 26, the solenoid valve is closed, stopping the flow of NO.\nFIG. 3 is a graph of the inspiratory gas flow 30 vs time for a conventional ventilator when the ventilator is in a second ventilating mode. When it is in this mode, the flow of gas is not constantly delivered for a predetermined duration but follows a particular curve during the inspiratory phase 32 and the patient then expires (expiratory phase 34). In the injection system proposed by Dubé et al., when the gas flow reaches a predetermined threshold level 36, the solenoid valve is open delivering NO to the patient. The line 38 illustrates the flow of NO over time. Again, it is to be noted that the scale is different for the flow of inspiratory gas and the flow 38 of NO. Indeed, line 38 illustrating the flow of NO is shown scaled up for illustrative purposes.\nSince the solenoid valve used by Dubé et al. is of the type fully open/fully closed, the flow of NO is constant when the valve is open. As can be seen from FIG. 4, the concentration of NO (line 39) is not constant over time during the inspiratory phases but varies inversely with the flow of gas since the flow of NO is constant. When the inspiratory Gas flow 30 falls below the threshold level 36, the solenoid valve is closed.\nA drawback of the injection system of Dubé et al. is that, in certain cases, the NO concentration is not constant during the inspiratory phase.\nCanadian patent application No. 2,133,516 filed on Oct. 3rd, 1994 and naming Bathe et al. as inventors describes a nitric oxide (NO) delivery system monitoring the inspiratory gas flow of a patient and controlling a proportional valve to allow a calculated flow of NO to enter the inspiratory gas flow. The delivery system calculates the flow of NO in order to maintain a constant, user programmable, NO concentration in the inspiratory gas.\nA drawback of the delivery system of Bathe et al. is that, while the delivery system may be programmed so that the concentration of NO in the inspiratory gas flow is constant, there are no provisions to modify the concentration of the NO during a particular inspiratory phase, or to program the variation of the concentration of NO over a number of successive inspiratory phases in view of gradually increasing or decreasing the concentration of NO supplied to the patient."}
{"text": "In coal and other kinds of mining by the Longwell technique, it is conventional for minerals to be removed by a single or double ended ranging shearer drum, which traverses the mineral face with a rotary cutting head carried by the, or each, ranging arm to follow the seam. Typically, each drum is provided with 50 or more locations where a cutting tool is required. A pick holder is welded in place at each location. Each pick holder supports a replaceable pick designed to engage the ground. In some constructions, each pick holder also contains a water sprayer to the rear of the pick for spraying the working end (i.e., the head) of the pick and the coal with water. In general, each pick comprises a pick shank, a securement mechanism to maintain the pick in the pick holder, a head, and a transition area between the head and the shank. The transition area often consists of a rear heel and a forward toe or shoulder.\nIn use, the shearer drum is rotated about its central axis. As the drum rotates the pick holders spin around with the drum so that picks engage the ground. The water sprayer within the pick holder sprays water on the pick and the coal to minimize dust and the risk of frictional ignitions.\nWhen the pick contacts the wall while the shearer drum rotates, the picks experience forces F as the pick breaks up the material to be excavated. The force F will at times be normal N to the tip 24 with respect to the material face such as along a line 1a normal and orthogonal to the frontal direction. Line 1a goes through the forward most impact point of the pick 10a to the center of rotation of the pick assembly 8a around the excavating equipment 4. In this application, a force that is along line 1a and only has a normal component N is referred to as a normal (or inward) force and a force F that is collinear or tangent T to the cutting path (i.e., orthogonal to line 1a and only has a tangential component) is referred to as a tangential (or rearward) force.\nAs a pick rotates around with the drum 6, the pick will experience a force F that will at times be primarily tangential (i.e., a force that extends perpendicular to a force that extends normally through the strike point of the tip to the center of rotation of the pick assembly around the drum and has an angle α of 0 degrees). Other times the pick will experience a force at an angle α from tangential T that has a tangential component and a normal component (i.e., a force between tangential T and normal N). As the pick continues to rotate the pick will experience a force that is primarily normal (i.e., a force that is primarily inward on the pick and has an angle α of 90 degrees from tangential T and that extends normally through the strike point of the tip 24a to the center of rotation of the pick assembly around the drum 6). The transition of forces between those that are predominantly inward or normal and those that are predominantly rearward or tangential causes the pick to rock within the pick holder. The cyclic rocking causes the pick to wear the pick holder prematurely as can be seen in FIG. 4. Premature wear on the pick holder causes the pick locking system to become ineffective and leads to picks being ejected from the pick holder during use. Typically, when a pick breaks or is ejected the pick will also break the water sprayer.\nDamaged holders must be cut from the drum and new holders welded in their place. Because of the risk of frictional ignitions and tight dark working areas, typically, shearer drums must be removed from the longwall (i.e., removed from service) and moved to a safe location for refurbishment, for example to the surface. Moving the shearer drum, cutting the welds between the shearer drums and the pick holder, and welding new pick holders in place is time consuming. Such refurbishment can be lengthy and expensive."}
{"text": "This invention relates to medical guidewires typically used by physicians to gain access to restricted regions of the body and over which therapeutic devices are passed for insertion to a site of interest. Specifically, the invention relates to an ultrasound imaging guidewire with a detachable imaging guidewire body and a stationary central core.\nMany surgeries involve the insertion of guidewires into a patient\"\"s body. The guidewire may be inserted into the digestive tract, urethra, blood vessels, heart chamber, a body cavity such as the abdominal cavity, or a hollow organ. Typically, an artery is the vessel of interest. The artery could be a relatively large peripheral vessel, a relatively small proximal coronary artery, or an artery of any size in between. The guidewire may include an imaging portion that permits close examination of the site of interest by means of ultrasonic waves. An ultrasonic imaging guidewire may permit the user to obtain 360 degree (i.e., cross-sectional) acoustic images of the vessel wall to, for example, determine the tissue morphology state of a site of interest, position a therapeutic device, monitor the progress of treatment or observe the site after treatment to determine the course of further treatment.\nOften, the guidewire must be Positioned at a predetermined site after passing through a complex network of blood vessels. Such placement may require a considerable amount of time. Furthermore, the difficulty and time required for guidewire placement increases with increasing vessel occlusion at later stages of disease. Thus, placement of the guidewire can be a time-consuming and difficult task.\nAccordingly, once the physician has taken the time to correctly place the guidewire, it is preferable to maintain the guidewire position. However, it is also desirable to obtain images of the diseased area which may require that the guidewire be axially translated to view the site of interest. Hence, after the physician places the guidewire, the physician needs to move the imaging guidewire back and forth to make a correct diagnosis of the lesion morphology. The problem with advancements and pullbacks of the imaging guidewire is that the physician may lose the correct placement of the guidewire, and have to spend additional time repositioning the guidewire. Thus, there currently exists a need to maintain guidewire positioning while permitting multi-position, real-time imaging.\nFurthermore, the back-and-forth movement of the guidewire may damage the patient\"\"s vessels. Therefore, there currently exists a need to provide safer guidewire imaging.\nA significant problem encountered by physicians is the proper positioning of stents. Stents are often used to prevent lumen closure following bypass surgery and to treat acute vessel closure after angioplasty. It is often extremely difficult for a physician to accurately determine the correct location to deploy a stent, particularly at a bifurcating vessel. Incorrect placement of a stent can lead to xe2x80x9cstent jailxe2x80x9d and is demonstrated in FIG. 3. As shown in FIG. 3, if the stent 100 is incorrectly placed at a bifurcating vessel location 102, the stent 100 may block the vessel 102 and the physician can no longer access that vessel 102. This is particularly dangerous if the vessel 102 becomes diseased, such as at 104, and access is needed for therapy. Thus, there currently exists a need for easier, multi-position, ultrasonic imaging of the site of interest to assist in accurate placement of a stent.\nThere also currently exists a need to provide improved imaging capabilities, without losing proper guidewire positioning, so as to efficiently locate the site of interest, to properly position therapeutic catheters such as an angioplasty balloon, and to observe continuously the site or sites of interest. There also exists a need to decrease the complexity and to save time associated with the ultrasonic imaging procedure.\nAccordingly, a general object of the present invention is to provide an apparatus and method for permitting multi-position, ultrasonic imaging without losing correct guidewire positioning.\nA further object of this invention is to provide a faster imaging guidewire procedure, and to eliminate the complexity associated with the ultrasonic imaging guidewire procedure.\nAnother object of this invention is to prevent harm to a patient\"\"s vessels by eliminating the back and forth movement of the guidewire tip.\nIn order to achieve the above objects of the present invention, an ultrasound imaging guidewire is provided with a connector to permit a static central core to be temporarily detached from an imaging guidewire body of a guidewire. A method is also provided to permit efficient and accurate imaging of the site of interest. The method includes the step of inserting a guidewire with an imaging guidewire body and a static central core into a patients body at a particular site of interest. Next, the imaging guidewire body is rotated at the site of interest to obtain acoustical images. Finally, the imaging guidewire body of the guidewire is axially translated to further obtain images of the site or sites of interest, without axially translating the static central core.\nAdditional objects, advantages, aspects and features of the present invention will further become apparent to persons skilled in the art from a study of the following description and drawings."}
{"text": "An example of prior art chip capacitors is shown in FIG. 1. The chip capacitor shown in FIG. 1 includes a solid-state tantalum capacitor element 2 with a cathode layer 4 disposed on its outer surface. An anode lead 6 is led out from one end surface of the capacitor element 2. A flat cathode terminal 8 is connected to the cathode layer 4 with an electrically conductive adhesive (not shown). Also, a flat anode terminal 10 is welded to the tip end of the anode lead 6. An encapsulation 12 is provided by transfer molding with epoxy resin. Outer end portions of the flat anode and cathode terminals 10 and 8 are bent to extend along the end surfaces of the encapsulation 12 and, then, further bent to extend along the bottom surface of the encapsulation 12.\nIt is seen that a large proportion of the cathode terminal 8 is within the encapsulation 12, and the proportion of the volume occupied by the cathode terminal 8 to the entire volume of the encapsulation 12 is large. Further, both the cathode terminal 8 and the anode terminal 10 include portions extending on the side surfaces of the encapsulation 12. Accordingly, the length of the capacitor is increased by the thickness of these portions. In mounting such chip capacitor on a printed circuit board, the side surfaces of the cathode and anode terminals 8 and 10 are connected to the board by solder 14. Accordingly, when a number of such chip capacitors are to be mounted on a board side by side, as shown in FIG. 2, the spacing between adjacent chip capacitors must be large enough to prevent short-circuiting of adjacent capacitors, which prevents dense packing of the capacitors. Recently, smallsized, portable electric and electronic devices, such as cellular phones, have been remarkably improved, and chip capacitors to be used in such devices are required to be down-sized. For down-sizing prior art chip capacitors like the ones described above, the volume occupied by the capacitor element 2 in the chip capacitor including the encapsulation 12 should be as small as possible, which sometimes prevents the chip capacitor from having desired capacitance.\nTherefore an object of the present invention is to provide a chip capacitor which makes high density packing possible, and can have desired capacitance, while being small in size."}
{"text": "A large and growing population of users is enjoying entertainment through the consumption of digital media items, such as music, movies, images, electronic books, and so on. The users employ various electronic devices to consume such media items. Among these electronic devices (referred to herein as user devices) are electronic book readers, cellular telephones, personal digital assistants (PDAs), portable media players, tablet computers, netbooks, laptops and the like. These electronic devices wirelessly communicate with a communications infrastructure to enable the consumption of the digital media items. In order to wirelessly communicate with other devices, these electronic devices include one or more antennas.\nThe conventional antenna usually has only one resonant mode in the lower frequency band and one resonant mode in the high-band. One resonant mode in the lower frequency band and one resonant mode in the high-band may be sufficient to cover the required frequency band in some scenarios, such as in 3G applications. 3G, or 3rd generation mobile telecommunication, is a generation of standards for mobile phones and mobile telecommunication services fulfilling the International Mobile Telecommunications-2000 (IMT-2000) specifications by the International Telecommunication Union."}
{"text": "1. Field of the Invention\nGenerally, the present disclosure relates to integrated circuits, and, more particularly, to the manufacture of field effect transistors in complex circuits including memory areas, for instance, in the form of a cache memory of a CPU.\n2. Description of the Related Art\nIntegrated circuits comprise a large number of circuit elements on a given chip area according to a specified circuit layout, wherein transistor elements represent one of the major semiconductor elements in the integrated circuits. Hence, the characteristics of the individual transistors significantly affect overall performance of the complete integrated circuit. Generally, a plurality of process technologies are currently practiced, wherein, for complex circuitry, such as microprocessors, storage chips, ASICs (application specific ICs) and the like, MOS technology is currently one of the most promising approaches due to the superior characteristics in view of operating speed and/or power consumption and/or cost efficiency. During the fabrication of complex integrated circuits using MOS technology, millions of transistors, i.e., N-channel transistors and/or P-channel transistors, are formed on a substrate including a crystalline semiconductor layer. A MOS transistor, irrespective of whether an N-channel transistor or a P-channel transistor is considered, comprises so-called PN junctions that are formed by an interface of highly doped drain and source regions with an inversely or weakly doped channel region disposed between the drain region and the source region. The conductivity of the channel region, i.e., the drive current capability of the conductive channel, is controlled by a gate electrode formed above the channel region and separated therefrom by a thin insulating layer. The conductivity of the channel region, upon formation of a conductive channel due to the application of an appropriate control voltage to the gate electrode, depends on the dopant concentration, the mobility of the majority charge carriers and, for a given extension of the channel region in the transistor width direction, on the distance between the source and drain regions, which is also referred to as channel length. Hence, in combination with the capability of rapidly creating a conductive channel below the insulating layer upon application of the control voltage to the gate electrode, the conductivity of the channel region substantially determines the performance of the MOS transistors. Thus, the latter aspect renders the reduction of the channel length, and associated therewith the reduction of the channel resistivity, a dominant design criterion for accomplishing an increase in the operating speed of the integrated circuits.\nOn the other hand, the drive current capability of the MOS transistors also depends on the transistor width, i.e., the extension of the transistor in a direction perpendicular to the current flow direction, so that the gate length and thus the channel length, in combination with the transistor width, are dominant geometric parameters which substantially determine the overall transistor characteristics in combination with “transistor internal” parameters, such as overall charge carrier mobility, threshold voltage, i.e., a voltage at which a conductive channel forms below the gate insulation layer upon applying a control signal to the gate electrode, and the like. On the basis of field effect transistors, such as N-channel transistors and/or P-channel transistors, more complex circuit components may be created, depending on the overall circuit layout. For instance, storage elements in the form of registers, static RAM (random access memory), may represent important components of complex logic circuitries. For example, during the operation of complex CPU cores, a large amount of data has to be temporarily stored and retrieved, wherein the operating speed and the capacity of the storage elements have a significant influence on the overall performance of the CPU. Depending on the memory hierarchy used in a complex integrated circuit, different types of memory elements are used. For instance, registers and static RAM cells are typically used in the CPU core due to their superior access time, while dynamic RAM elements are preferably used as working memory due to the increased bit density compared to registers or static RAM cells. Typically, a dynamic RAM cell comprises a storage capacitor and a single transistor wherein, however, a complex memory management system is required so as to periodically refresh the charge stored in the storage capacitors which may otherwise be lost due to unavoidable leakage currents. Although the bit density of dynamic RAM devices may be very high, a charge has to be transferred from and to the storage capacitors in combination with periodic refresh pulses, thereby rendering these devices less efficient in terms of speed and power consumption compared to static RAM cells. Thus, static RAM cells may be advantageously used as high speed memory with moderately high power consumption, thereby, however, requiring a plurality of transistor elements so as to allow the reliable storage of an information bit.\nFIG. 1a schematically illustrates a circuit diagram of a static RAM cell 150 in a configuration as may typically be used in modern integrated circuits. The cell 150 comprises a storage element 151, which may include two inversely coupled inverters 152A, 152B, each of which may include a couple of transistors 100B, 100C. For example, in a CMOS device, the transistors 100B, 100C may represent an N-channel transistor and a P-channel transistor, respectively, while, in other cases, transistors of the same conductivity type, such as N-channel transistors, may be used for both the transistors 100B and 100C. A corresponding arrangement of N-channel transistors for the upper transistors 100C is illustrated at the right-hand side of FIG. 1a. Moreover, respective pass or pass gate transistors 100A may typically be provided so as to allow a connection to the bit cell 151 for read and write operations, during which the pass transistors 100A may connect the bit cell 151 to corresponding bit lines (not shown), while the gate electrodes of the pass transistors 100A may represent word lines of the memory cell 150. Thus, as illustrated in FIG. 1a, six transistors may be required to store one bit of information, thereby providing a reduced bit density for the benefit of a moderately high operating speed of the memory cell 150, as previously explained. Depending on the overall design strategy, the memory cell 150 may require the various transistor elements 100A, 100B, 100C to have different characteristics with respect to drive current capability in order to provide reliable operational behavior during read and write operations. For example, in many design strategies, the transistor elements are provided with minimum transistor lengths, wherein the drive current capability of the transistors 100B, which may also be referred to as pull-down transistors, may be selected to be significantly higher compared to the drive current capability of the pass transistors 100A, which may be accomplished by appropriately adjusting the respective transistor width dimensions for the given desired minimum transistor length.\nFIG. 1b schematically illustrates a top view of a portion of the memory cell 150 as a hardware configuration in the form of a semiconductor device. As illustrated, the device 150 comprises a silicon-based semiconductor layer 102, in which an active region 103 is defined, for instance, by providing a respective isolation structure 104 that laterally encloses the active region 103, thereby defining the geometric shape and size of the transistors 100A, 100B. As illustrated, the transistors 100A, 100B may be formed in and above the same active region 103 since both transistors may have the same conductivity type and may be connected via a common node, as is for instance illustrated as nodes 153A, 153B in FIG. 1a. As previously explained, the transistors 100A, 100B, i.e., the pass transistor and the pull-down transistor, may have substantially the same length so that respective gate electrodes 106 may have substantially the same length 106L, whereas a transistor width 103B of the pull-down transistor 100B may be greater compared to a transistor width 103A of the pass transistor 100A, in order to establish the different current capabilities of these transistors.\nFIG. 1c schematically illustrates a cross-sectional view taken along the line 1c of FIG. 1b. As illustrated, the device 150 comprises a substrate 101 which may typically be provided in the form of a silicon substrate, possibly in combination with a buried insulating layer (not shown) if a silicon-on-insulator (SOI) is considered. Above the substrate 101 and a possible buried insulating layer, the semiconductor layer 102, in the form of a silicon layer, is provided, in which the isolation structure 104 (see FIG. 1b) may be formed according to the desired shape so as to define the active region 103 according to the configuration as shown in FIG. 1b. That is, the active region 103 has the width 103B in the transistor 100B and has the width 103A in the transistor 100A. In this respect, an active semiconductor region is to be understood as a semiconductor region having an appropriate dopant concentration and profile so as to form one or more transistor elements in and above the active region, which have the same conductivity type. For example, the active region 103 may be provided in the form of a lightly P-doped semiconductor material, for instance in the form of a P-well, when the semiconductor layer 102 may extend down to a depth that is significantly greater than the depth dimension of the transistors 100A, 100B, when the transistors 100A, 100B may represent N-channel transistors. Similarly, the active region 103 may represent a basically N-doped region when the transistors 100A, 100B represent P-channel transistors. Furthermore, in the manufacturing stage shown in FIG. 1c, the transistors 100A, 100B may comprise the gate electrode 106, for instance in the form of a polysilicon material, which is separated from a channel region 109 by a gate insulation layer 108. Furthermore, depending on the overall process strategy, a sidewall spacer structure 107 may be formed on sidewalls of the gate electrodes 106. Additionally, drain and source regions 110 may be formed in the active region 103 and may connect the transistors 100A, 100B. Typically, metal silicide regions 111 are provided in the gate electrode 106 and an upper portion of the drain and source regions 110 so as to reduce contact resistance of these areas.\nThe device 150 is typically formed on the basis of the following processes. First, the isolation structure 104 may be formed, for instance, as a shallow trench isolation by etching respective openings into the semiconductor layer 102 down to a specific depth, which may even extend to a buried insulating layer, if provided. Thereafter, the corresponding openings may be filled with an insulating material by deposition and oxidation processes, followed by a planarization such as chemical mechanical polishing (CMP) and the like. During the process sequence for the isolation structure 104, advanced lithography techniques may have to be used in order to form a corresponding etch mask, which substantially corresponds to the shape of the active region 103, which requires the definition of a moderately narrow trench to obtain the desired reduced width 103A of the transistor 100A. Thereafter, the basic doping in the active region 103 may be provided by performing respective implantation sequences, which may also include sophisticated implantation techniques for introducing dopants for defining the channel doping and the like. Next, the gate insulation layers 108 and the gate electrodes 106 may be formed by depositing, oxidizing and the like an appropriate material for the gate insulation layer 106, followed by the deposition of an appropriate gate electrode material, such as polysilicon. Subsequently, the material layers are patterned by using advanced lithography and etch techniques, during which the actual length 106L of the gate electrodes 106 may be adjusted, thereby requiring extremely advanced process techniques so as to obtain a gate length of approximately 50 nm and less. Next, a part of the drain and source regions 110 may be formed by implanting appropriate dopant species followed by the formation of the spacer structure 107, or at least a portion thereof, followed by a subsequent implantation process for defining the deep drain and source areas, wherein a corresponding implantation sequence may be repeated on the basis of an additional spacer structure if sophisticated lateral concentration profiles may be required in the drain and source regions 110. Thereafter, appropriate anneal processes may be performed to re-crystallize implantation-induced damage in the active region 103 and also to activate the dopant species in the drain and source areas 110. It should be appreciated that, for a reduced gate length in the above-defined range, the sophisticated geometric configuration of the active region 103 may result in process non-uniformities, for instance during the deposition and etching of a spacer material for forming the sidewall spacer 107. Typically, the spacer structure 107 is formed by depositing an appropriate material, for instance a silicon dioxide liner (not shown) followed by a silicon nitride material, which may subsequently be selectively etched with respect to the silicon dioxide liner on the basis of well-established anisotropic etch recipes. However, at areas indicated as 112 in FIG. 1b, irregularities may be observed which may even be increased due to respective non-uniformities created during previously performed lithography processes, such as the lithography process for patterning the gate electrodes 106 and the like. Consequently, the areas 112 may have a significant influence on the further processing of the device 150, which may finally result in non-predictable behavior of the transistor 100B and thus the overall memory cell 150. For example, during the further processing, the metal silicide regions 111 may be formed by depositing a refractory metal, such as nickel, cobalt and the like, which may then be treated to react with the underlying silicon material, wherein typically the isolation structure 103 and the spacer structure 107 may substantially suppress the creation of a highly conductive metal silicide. However, due to the previously generated irregularities, respective leakage paths or even short circuits may be created, thereby undesirably influencing the final drive current capability of the transistor 100B, which may result in a less stable and reliable operation of the memory cell 150, thereby significantly contributing to yield loss of sophisticated semiconductor devices including static RAM areas.\nThe present disclosure is directed to various methods and devices that may avoid, or at least reduce, the effects of one or more of the problems identified above."}
{"text": "1. Field of the Invention\nThe invention relates to a drive unit, in particular for an injection unit or an ejector of an injection molding machine.\n2. Description of Related Art\nRecently, one has provided injection molding machines with electric and hydraulic drives, wherein actuations at high speed are exerted by the electric drive with relatively low forces, while the hydraulic drive is particularly advantageous if high axial forces have to be applied with comparatively minor actuations.\nIn the case of a closing unit of a plastics injection molding machine, for instance, the drive unit moves a movable tool faceplate of the machine. In so doing, the drive unit has to fulfill two important, different objects. On the one hand, it is to move the tool faceplate as quickly as possible for closing and for opening the mould so as to keep the cycle time of the manufacturing of an injection-molded component as short as possible. On the other hand, it is to impact the tool faceplate with a high clamping force, so that the tool can be kept shut against the high inner pressure during injection molding. The drive unit therefore has to be configured such that it is adapted to perform actuations at high speed and to apply high forces with a comparatively minor stroke. Requirements of this kind are posed, except with a closing unit, also with the actuation of ejectors or the injection unit of an injection molding machine.\nDE 101 21 024 A1 (cf. in particular FIGS. 26, 34) of the Applicant discloses a drive unit that is adapted to fulfill the afore-mentioned requirements. This drive unit comprises a hydraulic force transmitting element, the smaller piston unit of which is actuated via an electrically actuated stroke spindle device for closing a tool. This smaller piston unit may consist of one single smaller piston, or of a plurality of small pistons. These confine, along with a cylinder or interface and one or several large pistons of the force transmitting element, a pressure chamber, wherein, by the moving of the small piston unit into the pressure chamber, a high pressure can be generated, which acts, via the large active surface of the large pistons (power pistons) on the movable tool faceplate which may then be kept shut with high force. During the quick closing of the tool with comparatively low force, the interface is indirectly connected with a spindle nut of the spindle device, so that the piston unit with smaller diameter, the power piston, and the interface are jointly shifted by the spindle device. For applying the high force, the interface is fixed at the frame of the injection molding machine, so that the further closing movement of the tool is determined by the moving of the smaller piston unit into the pressure chamber and the corresponding axial movement of the large piston of the force transmitting element.\nIn one embodiment described in DE 101 21 024 A1 (FIG. 34), the coupling of the cylinder to the stroke spindle device is performed hydraulically. To this end, a chamber confined by a section of the small piston unit and the cylinder is impacted with pressure from a high pressure storage means, so that the pressure medium incorporated in the chamber acts like rigid pulling mechanics and the cylinder participates in the closing stroke of the stroke spindle device and thus of the small piston unit.\nIn an embodiment illustrated in FIG. 26 of DE 101 21 024 A1, the small piston unit is, during rapid motion, connected with the large piston via an electromagnetic coupling. This large piston is in turn centered with respect to the cylinder by a prestressed centering spring arrangement. The prestressing of this centering spring arrangement is chosen such that the axial shifting of the small piston unit is, during rapid motion, transferred to the large piston via the coupling, and from there via the centering spring arrangement to the cylinder so as to take it along.\nIn both known solutions the force transmitting element is designed to be double-acting, so that, for tearing open the tool, a high tear-open force acts on the tool via the force transmitting element as the small piston unit moves in opening direction. This movement of the small piston unit in opening direction is performed during the application of the tear-open force against the force of a prestressed pressure spring.\nA disadvantage of the initially mentioned known construction (FIG. 34) is that, for applying the high pressure in the chamber during rapid motion, a comparatively complex circuitry with high pressure storage means and electrically controlled direction control valve is required, so that this circuitry variant is very expensive and also requires substantial construction space.\nIn the solution illustrated in FIG. 26 of DE 101 21 024 A1, the large piston has to be designed with a very large surface due to the integrated coupling, so that a compact solution cannot be realized with such a construction."}
{"text": "A semiconductor module has been extensively used in a power converter, such as an electrical system of a hybrid vehicle or an electrical vehicle. The semiconductor module configuring an energy-saving controller has a power semiconductor element for controlling a high current.\nThe heat generated by such a power semiconductor element when controlling a high current tends to increase especially as miniaturization and power boosting of the power semiconductor element advance. Therefore, a major problem is to cool a semiconductor module having a plurality of power semiconductor elements.\nA liquid cooling cooler has conventionally been used in such a semiconductor module. The power semiconductor elements need to be cooled efficiently in order to improve the cooling efficiency of the semiconductor module. The liquid cooling cooler is designed in various ways to improve its cooling efficiency thereof, by increasing the flow rate of its refrigerant, shaping heat radiating fins (cooling body) into a shape to provide good heat-transfer efficiency, and increasing heat-transfer efficiency of materials configuring the fins.\nIncreasing the flow rate of refrigerant supplied to the cooler or adopting the fin structures providing good heat-transfer efficiency can easily increase a pressure loss of the refrigerant inside the cooler. Especially in a cooler that uses a plurality of heat sinks to cool a large number of power semiconductor elements, the pressure loss of the refrigerant is significant in a passage structure where refrigerant passages are connected in series. To reduce the pressure loss of the refrigerant, it is ideal to construct a cooler in which its cooling efficiency can be enhanced with a low refrigerant flow rate. However, adopting a new fin material for improving the heat-transfer efficiency of the fin material configuring the cooler can lead to increases in costs of the entire cooler.\nIn a recent cooler, a refrigerant introducing passage for introducing a refrigerant and a refrigerant discharge passage for discharging the refrigerant are arranged parallel to each other, and a plurality of heat sinks are disposed therebetween in a refrigerant circulation direction so as to be substantially perpendicular to the abovementioned passages (for example, Japanese Patent Application Publication No. 2001-35981, Japanese Patent Application Publication No. 2007-12722, Japanese Patent Application Publication No. 2008-205371, Japanese Patent Application Publication No. 2008-251932, Japanese Patent Application Publication No. 2006-80211, Japanese Patent Application Publication No. 2009-231677, Japanese Patent Application Publication No. 2006-295178). In this case, the refrigerant can flow parallel between fins configuring each heat sink, increasing the cooling performance per pressure loss and reducing the pressure loss of the refrigerant inside the passages, as shown in Japanese Patent Application Publication No. 2006-80211.\nJapanese Patent Application Publication No. 2009-231677 describes a liquid cooling cooler in which the entire rear-side wall of the casing is smoothly inclined forward from the right-side wall toward the left-side wall and in which the cross-sectional area of the passage of the entrance header part decreases gradually from the cooling liquid entrance side toward the left-side wall (see the paragraphs [0024] and [0031] and FIG. 2). Japanese Patent Application Publication No. 2008-205371 describes a liquid cooling cooler in which the connection water paths for introducing and discharging a refrigerant are disposed on the same side surface of the module and in which each of the paths is disposed in a direction perpendicular to the fins without changing the cross-sectional areas thereof (see FIG. 1).\nJapanese Patent Application Publication No. 2006-295178 describes a heat sink apparatus for use in a computer electronic device and the like. In this heat sink apparatus, the shape of the inflow guide plate extending toward the plurality of passages is configured so as to be inclined into the shape of a curve of a convex surface toward the plurality of passages, as it moves away from the inflow port. In addition, the cross-sectional area of the inflow guide part becomes small gradually from the inflow port. Moreover, the shape of the inflow guide plate is same as that of the inflow guide plate (see the paragraph [0030] and FIG. 6).\nIn the conventional cooling technologies, however, a drift distribution in which the refrigerant drifts away occurs due to the shapes of the heat sinks and refrigerant passages, the method for disposing the heater elements, or the shapes of the refrigerant introduction/discharge ports, etc. Such drift distribution caused in the conventional coolers disturbs the balance of the cooling performance. Therefore, a uniform and stable cooling performance cannot be accomplished. Another problem is that only the temperatures of heat generated in the semiconductor elements disposed opposing the refrigerant discharge port increase significantly, reducing the lives of the elements or damaging the elements.\nAs in the coolers disclosed in Patent Documents 6 and 7 in which the cross-sectional areas of the entrance header parts decrease gradually in a direction in which the entrance header parts extend, their flow rate distributions are improving, but the increase in the temperatures of the sections near the refrigerant introduction ports cannot be prevented."}
{"text": "Arrows conventionally include fletchings mounted on their rear ends to provide flight stability. Usually three and sometimes four fletchings are mounted in a circumferentially spaced relationship about the rear end of the arrow shaft. In addition, two fletchings and a single fletching have also been utilized on arrow shafts in the past but have never received any significant commercial acceptance by archers.\nFeathers were the only type of fletchings conventionally utilized by archers until about five years or so ago when \"rubber\" fletchings gained acceptance. Actually, the designation \"rubber\" fletching is now somewhat of a misnomer since this type of fletching is presently made from synthetic plastic, although such plastic does have some rubber-like characteristics. Usually rubber fletchings are extruded with a mounting foot and a vane projecting from the foot, and the fletching is cut after the extrusion to the required length with the vane having the desired shape. In addition to feather and rubber fletchings, sheet plastic has also been previously utilized to make fletchings but has never received any significant commercial acceptance by archers.\nArrow spin or rotation is desirable to maintain flight stability and is usually achieved by mounting fletchings on the arrow shaft either at a slight angle with respect to the elongated axis thereof or in a helical configuration thereabout such that a screw action takes place during forward flight through the air. The consequent rotation stabilizes the arrow flight even when subjected to head, side, and tail winds that would otherwise significantly alter the flight trajectory. Such arrow rotation is particularly important with hunting arrows whose flat blade type points can tend to \"sail\" if there is not sufficient stabilizing rotation.\nPrior art arrow fletchings of the type described above are illustrated by U.S. Pat. Nos.: 2,193,397; 2,277,743; 2,525,332; 3,106,400; 3,539,187; 3,595,579; 3,749,403; 3,895,802; 4,003,576; and 4,088,323.\nMy prior U.S. Pat. Nos. 3,756,602 and 4,012,043 disclose arrow fletchings made by sharply bending sheet plastic to the desired shape. The arrow fletching disclosed by my U.S. Pat. No. 3,756,602 patent includes vanes that are spaced outwardly from the shaft an increasing extent in the forward direction so as to provide a construction that compensates for cross-winds by steering the arrow into the wind. The arrow fletching disclosed by my U.S. Pat. No. 4,012,043 patent includes vanes that each define a pocket of decreasing volume from front to rear to effect a pressure buildup that causes stabilizing rotation of the arrow during flight.\nFrictional drag generated by arrow fletchings during flight is affected by wind changes and thus alters the flight trajectory. Wind changes are a much greater problem with the longer distances involved in target shooting as compared to hunting, since any change in the frictional drag due to wind changes is effective over a greater period of time with the higher trajectory required for longer distances. Also, arrow fletchings heretofore have not had a construction capable of compensating for wind changes in order to maintain the desired flight trajectory."}
{"text": "1. Field of the Invention\nThe present invention relates to a procedure for the production of enantiomerically pure monomeric and dimeric C-10 non-acetal derivatives of natural trioxane artemisinin having high in vitro antimalarial, antiproliferative and antitumor activities. The present invention further relates to the formation of a novel trioxane aldehyde compound produced via a chemoselective C--C bond formation at the C-10 position upon reaction of artemisinin trioxane lactone with lithiothiazole or lithiobenzothiazole. This trioxane aldehyde may then be reacted with organolithium, Grignard, and phosphorus ylide nucleophiles exclusively via carbonyl addition.\n2. Description of the State of Art\nEach year approximately 200-300 million people experience a malarial illness and over 1 million individuals die. In patients with severe and complicated disease, the mortality rate is between 20 and 50%.\nPlasmodium is the genus of protozoan parasites which is responsible for all cases of malaria and Plasmodium falciparum is the species of parasite that is responsible for the vast majority of fatal malaria infections. Malaria has traditionally been treated with quinolines such as chloroquine, quinine, mefloquine, and primaquine and with antifolates such as sulfadoxine-pyrimethamine. Unfortunately, most P. falciparum strains have now become resistant to chloroquine, and some, such as those in Southeast Asia, have also developed resistance to mefloquine and halofantrine; multidrug resistance is developing in Africa also.\nThe endoperoxides are a promising class of antimalarial drugs which may meet the dual challenges posed by drug-resistant parasites and the rapid progression of malarial illness. The first generation endoperoxides include artemisinin and several synthetic derivatives, discussed in further detail below.\nArtemisia annua L., also known as qinghao or sweet wormwood, is a pervasive weed that has been used for many centuries in Chinese traditional medicine as a treatment for fever and malaria. Its earliest mention, for use in hemorrhoids, occurs in the Recipes for 52 Kinds of Diseases found in the Mawangdui Han dynasty tomb dating from 168 B.C. Nearly five hundred years later Ge Hong wrote the Zhou Hou Bei Ji Fang (Handbook of Prescriptions for Emergency Treatments) in which he advised that a water extract of qinghao was effective at reducing fevers. In 1596, Li Shizhen, the famous herbalist, wrote that chills and fever of malaria can be combated by qinghao preparations. Finally, in 1972, Chinese chemists isolated from the leafy portions of the plant the substance responsible for its reputed medicinal action. This crystalline compound, called qinghaosu, also referred to as QHS or artemisinin, is a sesquiterpene lactone with an internal peroxide linkage.\nArtemisinin is a member of the amorphane subgroup of cadinenes and has the following structure (I). ##STR2##\nArtemisinin or QHS was the subject of a 1979 study conducted by the Qinghaosu Antimalarial Coordinating Research Group involving the treatment of 2099 cases of malaria (P. vivax and P. falciparum in a ratio of about 3:1) with different dosage forms of QHS, leading to the clinical cure of all patients. See, Qinghaosu Antimalarial Coordinating Research Group, Chin. Med J, 92:811 (1979). Since that time artemisinin has been used successfully in many thousand malaria patients throughout the world including those infected with both chloroquine-sensitive and chloroquine-resistant strains of P. falciparum. Assay of artemisinin against P. falciparum in vitro revealed that its potency is comparable to that of chloroquine in two Hanian strains (Z. Ye, et al., J. Trad. Chin. Med., 3:95 (1983)) and of mefloquine in the Camp (chloroquine-susceptible) and Smith (chloroquine-resistant) strains, D. L. Klayman, et al., J. Nat. Prod., 47:715 (1984).\nAlthough artemisinin is effective at suppressing the parasitemias of P. vivax and P. falciparum, the problems encountered with recrudescence, and the compound's insolubility in water, led scientists to modify artemisinin chemically, a difficult task because of the chemical reactivity of the peroxide linkage which is believed to be an essential moiety for antimalarial activity.\nReduction of artemisinin in the presence of sodium borohydride results in the production of dihydroartemisinin (II-1) or DHQHS, (illustrated in structure II below), in which the lactone group is converted to a lactol (hemiacetal) function, with properties similar to artemisinin. Artemisinin in methanol is reduced with sodium borohydride to an equilibrium mixture of .alpha.- and .beta.-isomers of dihydroartemisinin. The yield under controlled conditions is 79% (artemisinin, 0.85M with NaBH.sub.4 6:34M, 7:5 equivalents in methanol, 12 L at 0-5.degree. C. for about 3 hours followed by quenching with acetic acid to neutrality at 0-5.degree. C. and dilution with water to precipitate dihydroartemisinin), A. Brossi, et al., Journal of Medicinal Chemistry, 31:645-650 (1988). Using dihydroartemisinin as a starting compound a large number of other derivatives, such as, ##STR3##\n1 R=H\n2 R=CH.sub.3 PA0 3 R=CH.sub.2 CH.sub.3 PA0 4 R=COCH.sub.2 CH.sub.2 COON.sub.a PA0 5 R=CH.sub.2 C.sub.6 H.sub.4 COOH PA0 6 R=CH.sub.2 C.sub.6 H.sub.4 COON.sub.a ##STR4##\nartemether (compound II-2), arteether (II-3), sodium artesunate (II-4), artelinic acid (II-5), sodium artelinate (II-6), dihydroartemisinin condensation by-product (II-7) and the olefinic compound, structure III, ##STR5##\nhave been produced.\nArtemether (II-2) is produced by reacting .beta.-dihydroartemisinin with boron trifluoride (BF.sub.3) etherate or HCl in methanol:benzene (1:2) at room temperature. A mixture of .beta.- and .alpha.-artemether (70:30) is obtained, from which the former is isolated by column chromatography and recrystallized from hexane or methanol, R. Haynes, Transactions of the Royal Society of Tropical Medicine and Hygiene, 88(1): S1/23-S1/26 (1994). For arteether (II-3), (Brossi, et al., 1988), the .alpha.-isomer is equilibrated (epimerized) to the .beta.-isomer in ethanol:benzene mixture containing BF.sub.3 etherate. Treatment of dihydroartemisinin with an unspecified dehydrating agent yields both the olefinic compound, (III), and the dihydroartemisinin condensation by-product (II-7), formed on addition of dihydroartemisinin to (III), M. Cao, et al., Chem. Abstr., 100:34720k (1984). Until recently, the secondary hydroxy group in dihydroartemisinin (II-1) provided the only site in an active artemisinin-related compound that had been used for derivatization. See B. Venugopalan, \"Synthesis of a Novel Ring Contracted Artemisinin Derivative,\" Bioorganic & Medicinal Chemistry Letters, 4(5):751-752 (1994).\nThe potency of various artemisinin-derivatives in comparison to artemisinin as a function of the concentration at which the parasitemia is 90 percent suppressed (SD.sub.90) was reported by D. L. Klayman, \"Qinghaosu (Artemisinin): An Antimalarial Drug from China,\" Science 228:1049-1055 (1985). Dr. Klayman reported that the olefinic compound III is inactive against P. berghei-infected mice, whereas the dihydroartemisinin condensation by-product (II-7) has an SD.sub.90 of 10 mg/Kg in P. berghei-infected mice. Thus, the dihydroartemisinin ether dimer proved to be less potent than artemisinin, which has an SD.sub.90 of 6.20 mg/Kg. Following, in order of their overall antimalarial efficacy, are the three types of derivatives of dihydroartemisinin (II-1) that have been produced: (artemisinin)&lt;ethers (II, R=alkyl)&lt;esters [II, R=C(=O)-alkyl or -aryl]&lt;carbonates [II,R=C(=O)O-alkyl or -aryl].\nOther rational design of structurally simpler analogs of artemisinin has led to the synthesis of various trioxanes, some of which possess excellent antimalarial activity. Posner, G. H., et al., reported the chemistry and biology of a series of new structurally simple, easily prepared, racemic 1,2,4-trioxanes (disclosed in U.S. Pat. No. 5,225,437 and incorporated herein by reference) that are tricyclic (lacking the lactone ring present in tetracyclic artemisinin I) and that are derivatives of trioxane alcohol IV ##STR6##\nhaving the relative stereochemistry shown above. Especially attractive features of trioxane alcohol IV are the following: (1) its straightforward and easy preparation from cheap and readily available starting materials, (2) its availability on gram scale, and (3) its easy one-step conversion, using standard chemical transformations, into alcohol derivatives such as esters and ethers, without destruction of the crucial trioxane framework. See, Posner, G. H., et al., J. Med. Chem., 35:2459-2467 (1992), incorporated herein by reference. The complete chemical synthesis of artemisinin and a variety of other derivatives is reviewed by Sharma, R. P., et al., Heterocycles, 32(8):1593-1638 (1991), and is incorporated herein by reference.\nMetabolic studies by Baker, et al., demonstrated that .beta.-arteether (II-3), one of the antimalarial derivatives discussed previously, was rapidly converted by rat liver microsomes into dihydroartemisinin (II-1). See Baker, J. K., et al., Biol. Mass Spect., 20: 609-628 (1991). This finding and the fact that the most effective artemisinin derivatives against malaria have been ethers or esters of (II-1) suggest that they were prodrugs for (II-1). The controlled slow formation of (II-1), however, is not desirable in view of its short half-life in plasma (less than two hours) and relatively high toxicity.\nThe successful synthesis of anticancer and antiviral drugs by replacing a carbon-nitrogen bond in nucleosides by a carbon--carbon bond (C-nucleosides) prompted the preparation of several 10-alkyldeoxoartemisinins, V, ##STR7##\nwherein R is 1-allyl, propyl, methyl, or ethyl. Typically, these syntheses involved five or six steps and the reported yields were only about 12 percent. See, Jung, M., et al., Synlett., 743-744 (1990); and Haynes, R. K., et al., Synlett., 481-484 (1992).\nHeterolytic cleavage of the peroxide O- 13 O bond via S.sub.N 2 attack of nucleophiles is well documented., see Adam, W. et aL, J. Am. Chem. Soc., 114:5591 (1992) and Razuvaev, G. A., et al., T. G. In Organic Peroxides; D. Swern, Ed., John Wiley & Sons, N.Y., 3:141-270, (1972). For example, tert-butyl ethers are conveniently prepared by Grignard nucleophilic attack on the O--O bond in tert-butyl peresters, see Lawesson, S.-O., et al., J. Am. Chem. Soc., 81:4230, (1959). Also, 3,3-disubstituted-1,2-dioxetanes react with organolithium reagents primarily via S.sub.N 2 O--O bond cleavage (with regioselective attack at the sterically less encumbered O atom) to form .beta.-hydroxy ethers (Adam, W., et al., Chem. Ber, 125:235, (1992)) and bicyclic endoperoxides likewise react with lithium and magnesium organometallics to produce O--O bond-cleaved hydroxy ethers. See, Schwaebe, M. K., et al., Tetrahedron Lett., 37:6635 (1996). When a dialkyl peroxide O---O bond is sterically hindered, then nucleophilic attack by a reactive organometallic reagent is made more difficult; an excellent example of this phenomenon leading to chemoselective nucleophilic addition of an organolithium reagent to the aldehyde carbonyl group in a peroxy aldehyde. See, Dussault, P., et al., T. J. Org. Chem., 58:5469 (1993) and Dussault, P., Synlett, 997 (1995). 1,2,4-Trioxanes in the artemisinin family undergo peroxide O--O bond cleavage when exposed to dimethylcopperlithium and to trityllithium; in these two cases, however, single-electron-reductive cleavage of the peroxide bond is likely occurring. See, Posner, G. H. et al., J. Am. Chem. Soc., 114:8328 (1992). Sodium borohydride chemoselectively reduces artemisinin (I) into its lactol (II-1), but more potent lithium aluminum hydride reduces both the lactone carbonyl group and the trioxane O--O bond. See, Wu, Y, et al, Youji Huaxue, 153, (1986); Chem. Abstr. 1986, 105, 191426n.\nBased on these published precedents, it seemed that it would be very difficult to find any reactive organometallic reagents that would add chemoselectively to the lactone carbonyl group (less electrophilic than an aldehyde) of trioxane lactone artemisinin (I) without also cleaving the trioxane O--O bond. In fact, exposing artemisinin to 1.2 equivalent of phenyllithium in THF at -78.degree. C. produced at least three major products (not characterized).\nOver the past thirty years only a few drugs isolated from higher plants have yielded clinical agents, the outstanding examples being vinblastine and vincristine from the Madagascan periwinkle, Catharanthus roseus, etoposide, the semi-synthetic lignan, from Mayapple Podophyllum peltatum and the diterpenoid taxol, commonly referred to as paclitaxel, from the Pacific yew, Taxus brevifolia. Of these agents, paclitaxel is the most exciting, recently receiving approval by the Food and Drug Administration for the treatment of refractory ovarian cancer. Since the isolation of artemisinin, there has been a concerted effort by investigators to study other therapeutic applications of artemisinin and its derivatives.\nNational Institutes of Health reported that artemisinin is inactive against P388 leukemia. See NCI Report on NSC 369397 (tested on Oct. 25, 1983). Later anticancer studies that have been conducted on cell line panels consisting of 60 lines organized into nine, disease-related subpanels including leukemia, non-small-cell lung cancer, colon, CNS, melanoma, ovarian, renal, prostate and breast cancers, further confirm that artemisinin displays very little anticancer activity. A series of artemisinin-related endoperoxides were tested for cytotoxicity to Ehrlich ascites tumor (EAT) cells using the microculture tetrazolum (MTT) assay, H. J. Woerdenbag, et al., \"Cytotoxicity of Artemisinin-Related Endoperoxides to Ehrlich Ascites Tumor Cells,\" Journal of Natural Products, 56(6):849-856 (1993). The MTT assay, used to test the artemisinin-related endoperoxides for cytotoxicity, is based on the metabolic reduction of soluble tetrazolium salts into insoluble colored formazan products by mitochondrial dehydrogenase activity of the tumor cells. As parameters for cytotoxicity, the IC.sub.50 and IC.sub.80 values, the drug concentrations causing respectively 50% and 80% growth inhibition of the tumor cells, were used. Artemisinin (I) had an IC.sub.50 value of 29.8 .mu.M. Derivatives of dihydroartemisinin (II-1) being developed as antimalarial drugs (artemether (II-2), arteether (III-3), sodium artesunate (II-4), artelinic acid (II-5) and sodium artelinate (II-6)), exhibited a somewhat more potent cytotoxicity. Their IC.sub.50 values ranged from 12.2 .mu.M to 19.9 .mu.M. The dihydroartemisinin condensation by-product dimer (II-7), disclosed previously by M. Cao, et al., 1984, was the most potent cytotoxic agent, its IC.sub.50 being 1.4 .mu.M. At this drug concentration the condensation by-product (II-7) is approximately twenty-two times more cytotoxic than artemisinin and sixty times more cytotoxic than dihydroartemisinin (II-1), the parent compound.\nWhile artemisinin and its related derivatives (II-6) discussed above demonstrated zero to slight antiproliferative and antitumor activity, it has been discovered that a class of artemisinin dimer compounds exhibits antiproliferative and antitumor activities that are, in vitro, equivalent to or greater than known antiproliferative and antitumor agents. See, U.S. Pat. No. 5,677,468 incorporated herein by reference. Unfortunately, while the in vitro results of these artemisinin compounds are encouraging these compounds do not appear to have significant antitumor activity on the treatment of tumor cells in mice.\nThere is still a need, therefore, to develop methods for the formation of hydrolytically stable C-10 carbon-substituted artemisinin compounds and structural analogs thereof having antimalarial, and antiproliferative and antitumor activities that are equivalent to or greater than those of known antimalarial, and antiproliferative and antitumor agents, respectively, wherein the method does not result in cleavage of the trioxane O--O bond."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication device incorporating the MAP (Manufacturing Automation Protocol), an international standard communication protocol that has been defined in ISO Standard ISO/DIS 9506-1, which is useable in a factory automation (FA) environment.\n2. Description of the Background Art\nIn an automated factory, a variety of devices are employed in the manufacturing operation and the devices are joined through a local communication network into a factory system. Since certain devices may be more suitable than others to perform desired manufacturing operations, often the devices used in the factory system will be manufactured by different vendors. Accordingly, each such FA device, whether a factory computer, robot, numerical control (NC) machine, programmable logic controller (PLC), process control equipment, or the like, will have a different type of microprocessor, use different computer languages and execute customized programs. It is desirable that the internal processing and operation of each device should have little effect on the way the devices interact in the factory system and, in particular, how they communicate with each other. In order to provide a common basis for communication, all of the devices in the system must use a common message structure (\"syntax\") and use a common set of messages or \"semantics\" (i.e., the naming of and access to remote variables, program loading, job management, error reporting and the like).\nThe Manufacturing Message Specification (MMS) has been adopted as an international standard that permits programs to be written for a variety of factory system devices on the basis of common semantics and syntax. The MMS is specified in two parts comprising the message services (semantics) and the protocol (syntax). The message services are grouped into functional units that relate to the kinds of functions that are performed when an application (a program that performs some desired job) at one user location interacts with the local communication network for purposes of communicating with a user at another (remote) location. A total of 86 message services may be grouped according to the functions of context management (e.g., Initiate, Conclude, Abort, Reject, Cancel), remote variable services (e.g., Read-data, Write, Define Named Variable, etc.), program services (Initiate Download Sequence for a program, Load Domain Content, etc.), diagnostics (Status, etc.), operator communication (Input and Output), coordination between applications (Define Semaphore, etc.), file services (File Open, File Read, etc.), event management (Define Event Condition, etc.), journal management (Read Journal, Write Journal, etc.) and job management/device control (e.g., Start-robot movement, Stop, Resume, etc.). A detailed description of the MMS standard appears in \"MMS Tutorial by John R. Tomlinson, System Integration Specialists Company, Inc. (1987).\nFIG. 4A is a block diagram illustrating the connection of two stations, each having corresponding applications and being interconnected by a local communication network, as they would appear in an automated factory environment. The application in station A \"at one end\" of the network communicates, via a MMS provider (shown as MMS), a logic link controller (LLC), a media access controller (MAC) and a modem at each station that is connected to a local network, with the application in station B \"at the other end\" of the network. In conventional MMS terminology, for such communication, station A is the \"Client\" and requests station B as the \"Server\" to perform some application specific operation; the Server responds with information resulting from the operation as it is performed. Typically, the Client is a controller station and the Server is a FA device.\nFIG. 4B is a block diagram illustrating the arrangement of a conventional communication device employing a PLC (Programmable Logic Controller) 1 as an example of an FA (Factory Automation) device. Ordinarily, the PLC has limited storage capability and relies on outside storage media (e.g., disk storage) to store pertinent programming and variables. Reliance on outside storage media has the disadvantage that when a power failure or OFF condition is encountered by the PLC, the relationship between the PLC and its external storage media must be redefined at power ON.\nIn FIG. 4B, the numeral 2 indicates a MAP interface unit, serving as a communication device and being connected between a MAP network 3 and the PLC 1 via a PLC-dedicated bus 4. The MAP interface 2 comprises an MMS protocol 5 whose communication object is a named variable, rather than an address. A PLC driver 7 for accessing the PLC 1, and a local manager 8 for carrying out management functions also are found in the MAP interface unit 2.\nFinally, interface 2 includes a VMD (Virtual Manufacturing Device) 6 for converting the MMS protocol 5 into a protocol reflecting the resources and functionality of the real FA device, e.g., PLC 1 in the preferred embodiment, and performing a process corresponding to each MMS service. The VMD, as an abstract representation of a Server showing its external behavior, comprises four conventional abstract elements including Executive function, Capabilities, Program Invocations and Domains. The latter are dynamic in nature and come into existence and are removed from the system either by MMS Services or by local action. The Domains comprise instructions and/or data which is dedicated to specific resources, such as the portion of the machine or robot that is controlled. Services are provided for a Client to manipulate Domains that are defined at the MMS Server, such as the Initiate Download Sequence and Upload Segment services.\nIn the standard MMS specification, the Domain management services comprise a Domain Object attribute, which specifies a VMD Object-specific name or Domain Name to uniquely identify the Domain within the VMD, and a List Of Capability attribute, which is a list of implementation specific parameters necessary to partition the resources of the VMD.\nThe PLC 1 is equipped with a computer interface 11. PLC 1 includes a symbolic address variable registration section 12, and is connected to the MAP network 3 via the dedicated bus 4 and MAP interface unit 2. A controller and multiple FA devices, each representing a different station, may be connected to the MAP network 3 for communication therebetween.\nThe VMD 6, as a \"virtual device\" that serves as an abstract model of the MMS server application, provides a consistent basis for defining the MMS services for all devices. In the present case, VMD 6 models the externally visible behavior of the PLC 1 and comprises applications that provide several MMS services and are represented as units, including Define Named Variable/Delete Named Variable means 61 for defining and deleting a named variable convertibly into a symbolic address variable specific to the PLC 1. Also included in VMD 6 is named variable accessing means 62 and a variable conversion table 63 wherein a named variable is registered (stored) in correspondence with a symbolic address variable specific to the FA device.\nFIG. 5 is a flowchart showing the operation of the MAP interface unit 2 acting as the communication device known in the art. The operation of the MAP interface unit 2 will now be described in reference to FIG. 5.\nReferring to FIG. 5, when a request for a Define Named Variable service is received from a station B at the other end (not shown) that is connected to the MAP network 3 in Step 201, the MMS protocol 5 activates the Define Named Variable/Delete Named Variable means 61 in the VMD 6 in Step 202. As a result, for example, the Define Named Variable/Delete Named Variable means 61 may register a named variable, e.g., \"DATA001,\" into the variable conversion table 63 in correspondence with a symbolic address \"D1\" according to the request of the other-end station B in Step 203. The named variable is related to a particular FA device, e.g., robot 1, as contrasted to robot 2 which may be represented by named variable \"DATA002\", and each FA device may be made by any of several vendors. Accordingly, the named variable is identified as having a relationship to a symbolic address, which ordinarily is vendor specific, e.g., Mitsubishi Electric Company of Japan has the standard address D1 and other unique standard addresses are assigned to other vendors. If the request is for a Delete Named Variable service, a corresponding named variable is deleted from the variable conversion table 63.\nWhen a request for a variable access service to the named variable \"DATA001\" is then received from the other-end station B in Step 204, the MMS protocol 5 activates the named variable accessing means 62 in the VMD 6 in Step 205, the named variable accessing means 62 converts the named variable \"DATA001\" into the symbolic address \"D1\" using the variable conversion table 63, and the VMD 6 accesses the symbolic address \"D1\" of the PLC 1 via the PLC driver 7 in Step 206. By using the table 63 which defines a named variable (e.g., DATA001) to be a vendor specific address (e.g., D1) programming is simplified and is useable for any of several devices from different vendors, since only a data call that is generic to the FA device at a given location (i.e., using DATA001) is used in the program to identify a desired operation, rather than a particular vendor address.\nSince the named variable is defined in a procedure as shown in Steps 201 to 203, i.e., when a request for the Define Named Variable service is received from the other-end station B (not shown), the MMS protocol 5 activates the Define Named Variable/Delete Named Variable means 61 in the VMD 6 to cause the Define Named Variable/Delete Named Variable means 61 to register the named variable into the variable conversion table 63 in response to the request of the other-end station B, registration cannot be made from other than the other-end station B. Accordingly, an application concerning a named variable to be registered for the other-end station B must be added for registration.\nSeveral other problems also are encountered in the conventional system design. For example, while a total of 86 services are set forth in the MMS, services which historically experience a low request level may not be provided. In fact, the actually provided services often comprise only about half of the total available services, due to the limited memory capacity in the PLC. For example, the other-end station B often is not provided with the Define Named Variable service or with the Delete Named Variable service. In the absence of these services, the table 63 cannot be utilized effectively, particularly when a power outage or OFF condition is encountered.\nMoreover, the known communication device arranged as described above does not allow a user-defined named variable for accessing an FA device to be registered from other than the other-end station. This requires an application for registering the named variable to be added to the other-end station for the purpose of registration."}
{"text": "Introduction to DRAs\nDielectric resonator antennas are resonant antenna devices that radiate or receive radio waves at a chosen frequency of transmission and reception, as used for example in mobile telecommunications. In general, a DRA consists of a volume of a dielectric material (the dielectric resonator) disposed on or close to a grounded substrate, with energy being transferred to and from the dielectric material by way of monopole probes inserted into the dielectric material or by way of monopole aperture feeds provided in the grounded substrate (an aperture feed is a discontinuity, generally rectangular in shape, although oval, oblong, trapezoidal ‘H’ shape, ‘<->’ shape, or butterfly/bow tie shapes and combinations of these shapes may also be appropriate, provided in the grounded substrate where this is covered by the dielectric material. The aperture feed may be excited by a strip feed in the form of a microstrip transmission line, grounded or ungrounded coplanar transmission line, triplate, slotline or the like which is located on a side of the grounded substrate remote from the dielectric material). Direct connection to and excitation by a microstrip transmission line is also possible. Alternatively, dipole probes may be inserted into the dielectric material, in which case a grounded substrate may not be required. By providing multiple feeds and exciting these sequentially or in various combinations, a continuously or incrementally steerable beam or beams may be formed, as discussed for example in the present applicant's co-pending U.S. patent application Ser. No. 09/431,548 and the publication by KINGSLEY, S. P. and O'KEEFE, S. G., “Beam steering and monopulse processing of probe-fed dielectric resonator antennas”, IEE Proceedings—Radar Sonar and Navigation, 146, 3, 121–125, 1999, the full contents of which are hereby incorporated into the present application by reference.\nThe resonant characteristics of a DRA depend, inter alia, upon the shape and size of the volume of dielectric material and also on the shape, size and position of the feeds thereto. It is to be appreciated that in a DRA, it is the dielectric material that resonates when excited by the feed, this being due to displacement currents generated in the dielectric material. This is to be contrasted with a dielectrically loaded antenna, in which a traditional conductive radiating element is encased in a dielectric material that modifies the resonance characteristics of the radiating element, but without displacement currents being generated in the dielectric material and without resonance of the dielectric material.\nDRAs may take various forms and can be made from several candidate materials including ceramic dielectrics.\nIntroduction to DRA Arrays\nSince the first systematic study of dielectric resonator antennas (DRAs) in 1983 [LONG, S. A., McALLISTER, M. W., and SHEN, L. C.: “The Resonant Cylindrical Dielectric Cavity Antenna”, IEEE Transactions on Antennas and Propagation, AP-31, 1983, pp 406–412], interest has grown in their radiation patterns because of their high radiation efficiency, good match to most commonly used transmission lines and small physical size [MONGIA, R. K. and BHARTIA, P.: “Dielectric Resonator Antennas—A Review and General Design Relations for Resonant Frequency and Bandwidth”, International Journal of Microwave and Millimetre-Wave Computer-Aided Engineering, 1994, 4, (3), pp 230–247].\nThe majority of configurations reported to date have used a slab of dielectric material mounted on a grounded substrate or ground plane excited by either a single aperture feed in the ground plane [ITTIPIBOON, A., MONGIA, R. K., ANTAR, Y. M. M., BHARTIA, P. and CUHACI, M: “Aperture Fed Rectangular and Triangular Dielectric Resonators for use as Magnetic Dipole Antennas”, Electronics Letters, 1993, 29, (23), pp 2001–2002] or by a single probe inserted into the dielectric material [McALLISTER, M. W., LONG, S. A. and CONWAY G. L.: “Rectangular Dielectric Resonator Antenna”, Electronics Letters, 1983, 19, (6), pp 218–219]. Direct excitation by a transmission line has also been reported by some authors [KRANENBURG, R. A. and LONG, S. A.: “Microstrip Transmission Line Excitation of Dielectric Resonator Antennas”, Electronics Letters, 1994, 24, (18), pp 1156–1157].\nThe concept of using a series of DRAs to build an antenna array has already been explored by several authors. For example, an array of two cylindrical single-feed DRAs has been demonstrated [CHOW, K. Y., LEUNG, K. W., LUK, K. M. AND YUNG, E. K. N.: “Cylindrical dielectric resonator antenna array”, Electronics Letters, 1995, 31, (18), pp 1536–1537] and then extended to a square matrix of four DRAs [LEUNG, K. W., LO, H. Y., LUK, K. M. AND YUNG, E. K. N.: “Two-dimensional cylindrical dielectric resonator antenna array”, Electronics Letters, 1998, 34, (13), pp 1283–1285]. A square matrix of four cross DRAs has also been investigated [PETOSA, A., ITTIPIBOON, A. AND CUHACI, M.: “Array of circular-polarized cross dielectric resonator antennas”, Electronics Letters, 1996, 32, (19), pp 1742–1743]. Long linear arrays of single-feed DRAs have also been investigated with feeding by either a dielectric waveguide [BIRAND, M. T. AND GELSTHORPE, R. V.: “Experimental millimetric array using dielectric radiators fed by means of dielectric waveguide”, Electronics Letters, 1983, 17, (18), pp 633–635] or a microstip [PETOSA, A., MONGIA, R. K., ITTIPIBOON, A. AND WIGHT, J. S.: “Design of microstrip-fed series array of dielectric resonator antennas”, Electronics Letters, 1995, 31, (16), pp 1306–1307]. This last research group has also found a method of improving the bandwidth of microstrip-fed DRA arrays [PETOSA, A., ITTIPIBOON, A., CUHACI, M. AND LAROSE, R.: “Bandwidth improvement for microstrip-fed series array of dielectric resonator antennas”, Electronics Letters, 1996, 32, (7), pp 608–609]. A study has also been made recently of different configurations that can be used to form cylindrical dielectric resonator antenna broadside arrays [WU, Z.; DAVIS, L. E. AND DROSSOS, G.: “Cylindrical dielectric resonator antenna arrays”, Proceedings of ICAP—11th International Conference on Antennas and Propagation, 2001, p. 668.]\nIt is important to note that the papers above have focused mainly on methods of feeding mechanisms for arrays of DRA elements and examining the benefits of such arrays for various applications. None of these publications has discussed the concept put forward in the present application, which is that of generating a specific DRA excitation mode in order to generate a specific far-field pattern that in turn enables a specific array geometry to be constructed.\nIntroduction to the Half-split DRA\nA problem with designing miniature dielectric resonator antennas for portable communications systems (e.g. mobile telephone handsets and the like) is that high dielectric materials must be used to make the antennas small enough to be physically compatible with the portable communications system. This in turn often leads to the antenna being too small in bandwidth. It is important therefore to identify DRA geometries and modes having low radiation quality factors and which are therefore inherently wide bandwidth radiating devices. It has been known for some time that the half-split cylindrical DRA is one such device see [JUNKER, G. P., KISHK, A. A. AND GLISSON A. W.: “Numerical analysis of dielectric resonator antennas excited in the quasi-TE modes”, Electronics Letters, 1993, 29, (21), pp 1810–1811] or [KAJFEZ, D. AND GUILLON, P.(Eds): “Dielectric resonators”, Artech House, Inc, Norwood, Mass., 1986.]. FIG. 1 of the present application shows the half-split DRA geometry and is taken from [KINGSLEY, S. P., O'KEEFE S. G. AND SAARIO S.: “Characteristics of half volume TE mode cylindrical dielectric resonator antennas”, to be published in IEEE Transactions on Antennas and Propagation, January 2002]. FIG. 1 shows a grounded conductive substrate 1 on which is disposed a half cylindrical dielectric resonator 2, with its rectangular surface 3 adjacent to the grounded substrate 1. The dielectric resonator 2 has a thickness d and a radius a, and is fed with a single probe 4 inserted into the rectangular surface 3 at a distance from a centre point of the surface 3. The resonator 2 also has a pair of semi-circular surfaces 5. The bandwidth of these half-split antennas has been the particular subject of a study [KISHK, A. A., JUNKER, G. P. AND GLISSON A. W.: “Study of broadband dielectric resonator antennas”, Published in Antenna applications Symposium, 1999, p. 45.] and bandwidths as high as 35% were reported for some configurations.\nUsing Half-split Cylindrical DRAs to Form an Array\nThe most common mode used for the half-split cylindrical DRA is the TE or quasi TE mode, which has the radiation patterns described in [KINGSLEY, S. P., O'KEEFE S. G. AND SAARIO S.: “Characteristics of half volume TE mode cylindrical dielectric resonator antennas”, to be published in IEEE Transactions on Antennas and Propagation, January 2002] or [JUNKER, G. P., KISHK, A. A. AND GLISSON A. W.: “Numerical analysis of dielectric resonator antennas excited in the quasi-TE modes”, Electronics Letters, 1993, 29, (21), pp 1810–1811]. In this mode, the direction of maximum radiation is along the long axis of the antenna. To form an antenna array from these elements, it is necessary to stack the elements 2 side by side with their long semi-circular faces 5 parallel to each other as shown in FIG. 2a. This gives minimum coupling between the elements 2—a requirement for good array design. This is a good way to form a horizontal array with vertical polarisation, but when the antenna array is turned vertically to from the type of array needed for mobile communications applications, for example, the array becomes horizontally polarised, as shown in FIG. 2b. Generally speaking, vertical polarisation is preferred to horizontal polarisation in many mobile communications applications as it gives better propagation at low elevation angles."}
{"text": "This invention relates to chainrings, and more particularly, to a solitary chainring for use with a conventional chain in a bicycle drivetrain system including a bicycle crank.\nBicycles and other chain-driven vehicles typically employ one or more chainrings and a set of rear hub mounted sprockets connected by a chain. Various mechanisms are used to maintain the chain on the chainring and sprockets. These mechanisms include chain guards, chain tensioners, chain catchers, derailleur configurations and so on.\nWhile riding a vehicle with a chain driven drivetrain, management of the chain and chainring engagement is particularly important to safe and effective propulsion of the bicycle. Keeping the chain engaged with the chainring can be difficult, which is especially true of geared bicycles which can experience severe changes in chain tension, and energetic motion of the chain, especially from riding over rough terrain.\nMoreover, the chainring in any bicycle can potentially touch the chain stay of the bicycle frame when the crank is in a position where high loads are applied by the rider, causing an elastic deformation of the bicycle frame and the crankset. This can lead to damage to the frame and chainring and cause other problems.\nThe invention provides an enhanced drive chain management, especially for a bicycle that can successfully and reliably be ridden over rough and challenging terrain."}
{"text": "The invention relates to the area of cancer diagnostics. More particularly, the invention relates to detection of the loss and or alteration of wild-type huBUB3 genes in tumor tissues.\nGenes and proteins involved in cell cycle regulation and apoptosis have been found to be important in the development of cancers. There is a continuing need in the art for identification of components of cells which control the cell cycle and apoptosis. These components can be used both diagnostically and therapeutically to identify and detect neoplasms as well as to treat them.\nIt is an object of the present invention to provide methods and tools for diagnosing and treating neoplasia. These and other objects of the invention are provided by one or more of the embodiments which are described below.\nOne embodiment of the invention is an isolated and purified huBUB3 protein having an amino acid sequence which is at least 85% identical to SEQ ID NO:2. Percent identity is determined using a Smith-Waterman homology search algorithm using an affine gap search with a gap open penalty of 12 and a gap extension penalty of 1.\nAnother embodiment of the invention is an isolated and purified polypeptide comprising at least 8 contiguous amino acids as shown in SEQ ID NO:2.\nEven another embodiment of the invention is a huBUB3 fusion protein comprising a first protein segment and a second protein segment fused together by means of a peptide bond. The first protein segment consists of at least 8 contiguous amino acids of a huBUB3 protein as shown in SEQ ID NO:2.\nStill another embodiment of the invention is a preparation of antibodies which specifically bind to a huBUB3 protein having an amino acid sequence as shown in SEQ ID NO:2.\nA further embodiment of the invention is a cDNA molecule which encodes a huBUB3 protein having an amino acid sequence which is at least 85% identical to SEQ ID NO:2. Percent identity is determined using a Smith-Waterman homology search algorithm using an affine gap search with a gap open penalty of 12 and a gap extension penalty of 1.\nYet another embodiment of the invention is a cDNA molecule which encodes at least 8 contiguous amino acids of SEQ ID NO:2.\nAnother embodiment of the invention is a cDNA molecule comprising at least 12 contiguous nucleotides of SEQ ID NO:1.\nStill another embodiment of the invention is a cDNA molecule which is at least 85% identical to the nucleotide sequence shown in SEQ ID NO:1. Percent identity is determined using a Smith-Waterman homology search algorithm using an affine gap search with a gap open penalty of 12 and a gap extension penalty of 1.\nEven another embodiment of the invention is an isolated and purified subgenomic polynucleotide comprising a nucleotide sequence which hybridizes to SEQ ID NO:1 after washing with 0.2xc3x97SSC at 65xc2x0 C. The nucleotide sequence encodes a huBUB3 protein having the amino acid sequence of SEQ ID NO:2.\nYet another embodiment of the invention is a construct comprising a promoter and a polynucleotide segment encoding at least 8 contiguous amino acids of a huBUB3 protein as shown in SEQ ID NO:2. The polynucleotide segment is located downstream from the promoter. Transcription of the polynucleotide segment initiates at the promoter.\nEven another embodiment of the invention is a host cell comprising a construct which comprises a promoter and a polynucleotide segment encoding at least 8 contiguous amino acids of a huBUB3 protein having an amino acid sequence as shown in SEQ ID NO:2.\nA further embodiment of the invention is a recombinant host cell comprising a new transcription initiation unit. The new transcription initiation unit comprises in 5xe2x80x2 to 3xe2x80x2 order an exogenous regulatory sequence, an exogenous exon, and a splice donor site. The new transcription initiation unit is located upstream of a coding sequence of a huBUB3 gene as shown in SEQ ID NO:1. The exogenous regulatory sequence controls transcription of the coding sequence of the huBUB3 gene.\nStill another embodiment of the invention is a pair of single stranded DNA primers. The set allows synthesis of all or part of a huBUB3 gene coding sequence.\nYet another embodiment of the invention is a nucleic acid probe complementary to a wild-type huBUB3 gene as shown in SEQ ID NO:1.\nEven another embodiment of the invention is a method of diagnosing a neoplastic tissue of a human. Loss of a wild-type huBUB3 gene or an expression product of the wild-type huBUB3 gene from a tissue suspected of being neoplastic is detected. The wild-type huBUB3 gene has the coding sequence shown in SEQ ID NO:1. The loss indicates neoplasia of the tissue.\nAnother embodiment of the invention is a method of identifying a neoplastic tissue of a human. Expression of a first huBUB3 gene in a first tissue of a human suspected of being neoplastic is compared with expression of a second huBUB3 gene in a second tissue of the human which is normal. The second huBUB3 gene has the coding sequence shown in SEQ ID NO:1. Decreased expression of the first huBUB3 gene relative to the second huBUB3 gene identifies the first tissue as being neoplastic.\nStill another embodiment of the invention is a method to aid in the diagnosis or prognosis of neoplasia in a human. A first huBUB3 gene, mRNA, or protein in a first tissue of a human suspected of being neoplastic is compared with a second huBUB3 gene, mRNA, or protein in a second tissue of a human which is normal. The second huBUB3 gene has the coding sequence shown in SEQ ID NO:1. A difference between the first and second huBUB3 genes, mRNAs, or proteins indicates the presence of neoplastic cells in the first tissue.\nEven another embodiment of the invention is a method to aid in detecting a genetic predisposition to neoplasia in a human. A huBUB3 gene, mRNA, or protein in the fetal tissue of a human is compared with a wild-type huBUB3 gene, mRNA, or protein. The wild-type huBUB3 gene has the coding sequence shown in SEQ ID NO:1. A difference between the huBUB3 gene, mRNA, or protein in the fetal tissue of the human and the wild-type huBUB3 gene, mRNA, or protein indicates a genetic predisposition to neoplasia in the human.\nYet another embodiment of the invention is a method of screening test compounds for the ability to interfere with the binding of a huBUB3 protein to a huBUB1 protein. A test compound is contacted with at least a huBUB3-binding domain of a huBUB1 protein as shown in SEQ ID NO:4 and at least a huBUB1-binding domain of a huBUB3 protein as shown in SEQ ID NO:2. The huBUB3-binding domain binds to the huBUB1-binding domain in the absence of the test compound. The amount of the huBUB1-binding domain which is bound or unbound to the huBUB3-binding domain or the amount of the huBUB3-binding domain which is bound or unbound to the huBUB1-binding domain in the presence of the test compound is determined. A test compound which decreases the amount of bound huBUB1- or huBUB3-binding domains or which increases the amount of unbound huBUB1- and huBUB3-binding domains is a potential inducer of mitosis or cell cycle progression.\nEven another embodiment of the invention is a method of screening test compounds for the ability to interfere with the binding of a huBUB1 protein to a huBUB3 protein. A cell is with a test compound. The cell comprises a first fusion protein, a second fusion protein, and a reporter gene. The first fusion protein comprises (1) at least a huBUB1-binding domain of a huBUB3 protein as shown in SEQ ID NO:2 and (2) either a DNA binding domain or a transcriptional activating domain. The second fusion protein comprises at least a huBUB3-binding domain of a huBUB1 protein as shown in SEQ ID NO:4. The huBUB1-binding domain binds to the huBUB3-binding domain. If the first fusion protein comprises a DNA binding domain, then the second fusion protein comprises a transcriptional activating domain. If the first fusion protein comprises a transcriptional activating domain, then the second fusion protein comprises a DNA binding domain. The interaction of the first and second fusion proteins reconstitutes a sequence-specific transcription activating factor. The reporter gene comprises a DNA sequence to which the DNA binding domain specifically binds. Expression of the reporter gene is measured. A test compound which decreases the expression of the reporter gene is a potential inducer of mitosis or cell cycle progression.\nAnother embodiment of the invention is a method of identifying compounds which interfere with the binding of a huBUB3 protein to a huBUB1 protein. A cell which comprises three recombinant DNA constructs is provided. A first construct encodes a first polypeptide fused to a sequence-specific DNA-binding domain, a second construct encodes a second polypeptide fused to a transcriptional activation domain, and a third construct comprises a reporter gene downstream from a DNA element which is recognized by the sequence-specific DNA-binding domain. The first polypeptide comprises a huBUB1-binding domain of a huBUB3 protein as shown in SEQ ID NO:2 and the second polypeptide comprises a huBUB3-binding domain of a huBUB1 protein as shown in SEQ ID NO:4 or the first polypeptide comprises a huBUB3-binding domain of a huBUB1 protein as shown in SEQ ID NO:4 and the second polypeptide comprises a huBUB1-binding domain of a huBUB3 protein as shown in SEQ ID NO:2. The cell is contacted with a test compound. Expression of the reporter gene in the presence of the test compound is determined. A test compound which decreases expression of the reporter gene is identified as a candidate therapeutic agent.\nYet another embodiment of the invention is a cell which comprises three recombinant DNA constructs. A first construct encodes a first polypeptide fused to a sequence-specific DNA-binding domain, a second construct encodes a second polypeptide fused to a transcriptional activation domain, and a third construct comprises a reporter gene downstream from a DNA element which is recognized by the sequence-specific DNA-binding domain. The first polypeptide comprises a huBUB1-binding domain of a huBUB3 protein as shown in SEQ ID NO:2 and the second polypeptide comprises a huBUB3-binding domain of a huBUB1 protein as shown in SEQ ID NO:4, or the first polypeptide comprises a huBUB3-binding domain of a huBUB1 protein as shown in SEQ ID NO:4 and the second polypeptide comprises a huBUB1-binding domain of a huBUB3 protein as shown in SEQ ID NO:2.\nEven another embodiment of the invention is a method of determining the quantity of huBUB1 which binds to huBUB3, or of huBUB3 which binds to huBUB1. A first protein and a second protein are contacted. If the first protein is huBUB3 then the second protein is huBUB1 and if the first protein is huBUB1 the second protein is huBUB3. The quantity of the first protein which is bound to the second protein is determined.\nStill another embodiment of the invention is a method for identifying compounds which decrease the kinase activity of a huBUB1-huBUB3 complex. A huBUB1-huBUB3 complex is contacted with a test compound. The kinase activity of the huBUB1-huBUB3 complex is determined. A compound which decreases kinase activity of the huBUB1-huBUB3 complex is identified as a candidate therapeutic agent.\nThe present invention thus provides the art with the sequence of the human BUB3 gene and protein. This information allows highly specific assays to be done to assess the neoplastic status of a particular tumor tissue."}
{"text": "This invention relates to a press-contacting terminal.\nFIGS. 4 and 5 show known related press-contacting terminals (see, for example, JP-A-2003-217700 (Pages 2 and 4 to 5, FIGS. 4 and 15)).\nThe press-contacting terminal 1, shown in FIG. 4, is formed by stamping and bending a single electrically-conductive metal sheet. This press-contacting terminal 1 includes a terminal connection portion 2 of a generally rectangular tubular shape into which a mating terminal can be inserted to be electrically connected thereto, and a wire connection portion 3 extending from the terminal connection portion 2 in a direction of connecting of the mating terminal to the terminal 1.\nThe wire connection portion 3 includes a base plate portion 4, a pair of side plate portions 5 formed on and projecting perpendicularly respectively from widthwise-opposite side edges of the base plate portion 4, and press-contacting portions 6 for cutting an insulating sheath of a sheathed wire to be brought into contact with an internal conductor of the sheathed wire. Each press-contacting portion 6 includes a pair of press-contacting blades 8 which are formed by stamping predetermined opposed portions of the pair of side plate portions 5 and then by bending these stamped-out portions inwardly. A slot 9 is formed between the pair of press-contacting blades 8. The two press-contacting blade portions 6 are arranged in the mating terminal-connecting direction.\nThe sheathed wire is press-fitted into a space between the pair of side plate portions 5 from the upper side, and the insulating sheath of the sheathed wire is cut by upper end portions of the press-contacting blades 8, and then the internal conductor of the sheathed wire is inserted into the slots 9, so that the sheathed wire is press-contacted with the press-contacting portions 6, and therefore is electrically connected to the press-contacting terminal 1.\nLike the above-mentioned press-contacting terminal 1, the press-contacting terminal 11 of FIG. 11 is formed by stamping and bending a single electrically-conductive metal sheet. This press-contacting terminal 11 includes a terminal connection portion 12 of a generally rectangular tubular shape into which a mating terminal can be inserted to be electrically connected thereto, and a wire connection portion 13 extending from the terminal connection portion 12 in a direction of connecting of the mating terminal to the terminal 11.\nTwo press-contacting piece portions 17a and 17b, each having a pair of press-contacting blades 18 cooperating with each other to define a slot 19 therebetween, are formed on the wire connection portion 13, and are arranged in the mating terminal-connecting direction. The terminal connection portion 12 is formed on a front side portion of the press-contacting terminal 11, and the rear press-contacting piece portion 17a is formed by bending a distal end portion of a base plate portion 14, extending from a base plate portion of the terminal connection portion 12, at a right angle. A strip-like portion 15 extends from a top plate portion of the terminal connection portion 12, and is bent to be superposed on the base plate portion 14 in a direction of the height, and the front press-contacting piece 17b is formed by bending a distal end portion of the strip-like portion 15 at a right angle.\nA sheathed wire is pressed down from the upper side of the press-contacting blades 17a and 17b, and an insulating sheath of the sheathed wire is cut by upper end portions of the press-contacting blades 18, and then an internal conductor of the sheathed wire is inserted into the slots 19, so that the sheathed wire is press-contacted with the press-contacting piece portions 17a and 17b, and therefore is electrically connected to the press-contacting terminal 11.\nThe above-mentioned press-contacting terminals 1 and 11 are mounted within a connector housing of a connector, and with an increasing demand for a compact design of recent connectors, it has now been required to achieve a low-height design of the press-contacting terminals (that is, the reduction of the size of the press-contacting terminals in the direction of their height) and a narrow-pitch arrangement of the terminals (that is, the reduction of the size of the presscontacting terminals in the direction of their width). However, the related press-contacting terminals, shown in FIGS. 4 and 5, have the following problems which are to be solved in order to meet the above requirements.\nIn the press-contacting terminal 1 of FIG. 4, the pair of press-contacting blades 8 of each press-contacting portion 6 are formed by bending the stamped-out portions of the pair of side plate portions 5, and are disposed between the pair of side plate portions 5. The sheathed wire to be press-contacted with the terminal is press-fitted into the space between the pair of side plate portions 5. Therefore, the width of the press-contacting terminal 1 is at least the sum of the width of the sheathed wire and the thicknesses of the pair of side plate portions 5 (that is, a thickness twice larger than the thickness of the metal sheet forming the press-contacting terminal 1). Therefore, the reduction of the press-contacting terminal 1 in the direction of the width thereof is limited, and this is disadvantageous with respect to the achievement of the narrow-pitch arrangement.\nOn the other hand, in the press-contacting terminal 11 of FIG. 5, the front press-contacting piece portion 17b is formed by bending the distal end portion of the strip-like portion 15 (superposed on the base plate portion 14 in the direction of the height) at a right angle. Therefore, the thickness of the strip-like portion 15 (that is, the thickness of the metal sheet forming the press-contacting terminal 11) is interposed between the press-contacting piece portion 17b and the base plate portion 14 in the direction of the height. This is disadvantageous with respect to the reduction of the press-contacting terminal 11 in the direction of the height. And besides, the wire connection portion 13 of the press-contacting terminal 11 has a flat plate-like shape, and therefore there is a fear that its strength is reduced."}
{"text": "1. Field of the Invention\nThe present invention relates to an audio coding device and a method thereof by which an input audio signal is coded according to so-called transform coding and an obtained code string is transferred or recorded onto a recording medium, and also relates to an audio decoding device and a method thereof by which a code string transferred or red from a recording medium is decoded to obtain an output audio signal.\n2. Description of the Related Art\nThere has been a known method in which spectrums obtained by performing time-frequency transform on an input audio signal are subjected to normalization/quantization and differential frequency spectrums as quantization errors are subjected again to normalization/quantization (see Patent Documents 1 and 2: Japanese Patent Publications No. 3227945 and No. 3227948). Quantization accuracy of the audio coding device can be improved by this method, and scalability can be realized to fit performance and use environment of the audio decoding device."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus for adjusting the relative humidity (hereinafter referred to simply as \"humidity\") of gas to a fixed value.\n2. Description of the Prior Art\nWithin a tightly closed space, retention of a gas at a fixed humidity can be effected rather easily by use of a saturated aqueous solution. In retaining a gas within a container at a fixed humidity by continuously introducing into the container gas having a substantially constant humidity, however, a large volume of the gas must be continuously supplied to the container. This supply of the gas proves to be quite difficult.\nIt has been customary to control the humidity of a given gas by means of a constant temperature vessel which incorporates a humidifier and a dehumidifier. In this constant temperature vessel, the gas is automatically maintained at a constant humidity by having the humidifier and the dehumidifier properly started or stopped alternately by means of a sensor and a controller. This method of humidity control, however, requires adoption of a prohibitively expensive apparatus.\nAn expeditious measure resorted to for the regulation of the humidity of a gas involves the use of two containers, one filled with sulfuric acid and the other with water. A given gas whose humidity is desired to be regulated is divided into two streams, one stream to be passed through the container filled with sulfuric acid and the other stream through the container filled with water, whereafter the streams from the two containers are combined and mixed. It may appear that it would be possible to freely vary the humidity of the mixed gas by suitably selecting the ratio at which the gas is divided into the two streams prior to passage through the two respective containers so that the humidity of the gas could be readily regulated automatically. In actuality, however, it is extremely difficult to have the gas introduced accurately at prescribed ratios into the two containers solely by manual handling of cocks adapted to adjust the apertures in the respective feed pipes. The manual handling of these cocks is quite susceptible of error. Further, the humidity cannot easily be kept at a constant value because the pressure of the saturated vapor is variable with the ambient temperature. This fact constitutes itself another drawback for the method under discussion.\nOne object of this invention is to provide a method for stably and easily adjusting the humidity of gas to a constant value.\nAnother object of this invention is to provide an apparatus for easily adjusting the humidity of gas to a constant value."}
{"text": "1. Field of the Invention\nThis invention relates to a form assembly, more particularly to a modular form assembly which has a plurality of vertical channel pieces that can be coupled detachably side by side to one another to constitute a desired dimension of common form plane for forming a concrete structure.\n2. Description of the Related Art\nThe improvement of this invention is directed to a conventional form assembly, as shown in FIG. 1, for forming a concrete structure. The conventional form assembly, which has been disclosed in U.S. Pat. No. 4,957,272, includes vertical form plates (A1), vertical backing frames (A2), and horizontal reinforcement channel members (A3) which are coupled detachably to one another and which are then attached detachably to upper and lower edges of the form plates (A1). The backing frames (A2) are mounted detachably between the horizontal reinforcement channel members (A3) behind the form plates (A1). In order to further reinforce the form plates (A1), the conventional form assembly, when in practical application, further includes several horizontal reinforcement rods (A4), as shown in FIGS. 2 and 3, which are mounted threadably to the backing frames (A2) to press against the rear faces of the form plates (A1) in order to increase the contact areas between the form plates (A1) and the backing frames (A2), thereby preventing deformation of the form plates (A1) during use.\nReferring again to FIG. 1, the form plates (A1) have positioning holes (A5) formed therethrough at predetermined positions where, when any two form plates (A1) are located face to face, the positioning holes (A5) can be aligned with each other. In order to maintain two opposite form plates (A1) at a predetermined space into which concrete is poured, the conventional form assembly further includes several positioning bars (not shown), each of which having its two end portions mounted respectively within the aligned positioning holes (A5) of the opposite form plates (A1). In this way, the predetermined space between the opposite form plates (A1) can be maintained.\nThe drawbacks of the conventional form assembly are as follows:\n1. A relatively large number of parts is required, such as the form plates (A1), the backing frames (A2), the reinforcement channel members (A3) and the reinforcement rods (A4), to constitute the conventional form assembly. Thus, it is quite complicated and difficult to assemble or disassemble these parts when in use.\n2. Because the dimensions of the form plates (A1) are relatively large, the form plates (A1) can only be applied for forming wider walls, such as partition walls (B1), as shown in FIG. 4. When applying the form plate (A1) to form narrower walls, such as the narrow walls (B3) at the aisles (B2) or at the vent passageway (B4), the form plate (A1) has to be cut to fit the dimensions of the narrow walls (B3). This may result in waste of material. Therefore, bricks are still employed to constitute the narrow walls (B3).\n3. To form variable dimensions of partition walls (B1) for use in different floors of a building, some of the form plates (A1) have to be cut so as to allow the assembly of the modified and original form plates (A1) to obtain the desired partition walls (B1). This results in waste of time and material. In addition, due to the change in the dimensions of the form plates (A1), alignment of the positioning holes (A5) at the opposite form plates (A1) cannot be maintained. Accordingly, several additional positioning holes (A5) have to be provided again in some of the form plates (A1) so as to mount the positioning bars. At the same time, the original positioning holes (A5) have to be filled. Thus, it is quite inconvenient to assemble the modified and original form plates (A1)."}
{"text": "It is often desirable to remove acid gases, such as, for example, CO.sub.2, H.sub.2 S, SO.sub.2, CS.sub.2, HCN, COS, and sulfur derivatives of C.sub.1 to C.sub.8 hydrocarbons, from gas streams. Gas streams from which these acid gases must be removed can be from many sources. One common source of such gas streams is from natural gas wells. The gas removed from natural gas wells is often rich in methane and other combustible gases, but contains concentrations of acid gases such as H.sub.2 S, CO.sub.2 and the other acid gases described above. High concentrations of H.sub.2 S inhibit pipe line shipment of the natural gas because of environmental considerations and government regulation. High concentrations of CO.sub.2 in natural gas reduce the heating value of the gas because CO.sub.2 is not combustible. Mercpatans, i.e., sulfur derivatives of C.sub.1 to C.sub.8 hydrocarbons, have an offensive odor and are corrosive.\nThe removal of mercaptans can be particularly difficult. One process proposed for the removal of mercaptans from a gas stream is described in U.S. Pat. No. 3, 716,620, issued Feb. 13, 1973. The process includes the step of contacting a gas containing hydrogen sulfide or a mercaptan with a solution of iodine in an organic solvent, e.g., an ether of a polyalkylene glycol, and an amine. The presence of iodine in processes such as described in the above-referenced patent is generally undesirable because the iodine must be regenerated in an oxidation process which increases the complexity and adds cost to the overall acid gas removal process.\nAccordingly, processes and absorption solvents are desired for the removal of mercaptans from gas streams by absorption which do not require the presence of iodine or suffer the disadvantages described above."}
{"text": "In conventional pulse width modulation systems a ramp is intersected by two variable references to accomplish dual edge modulation (DEM), by a fixed and a variable reference for trailing edge modulation (TEM), and by a variable and a fixed reference for leading edge modulation (LEM). Generally such systems operate satisfactorily but inaccuracies and unreliability can occur at the extremes: when zero width pulses and maximum width pulses are generated. For example, when maximum width pulses are requested the full width available from the ramp cannot be utilized because there must be some overshoot of the ramp to intersect the reference. When zero width pulses are requested it is likewise difficult to rely on the ramp instantaneously responding to the references to emit absolutely no pulse: some small spikes may get through. Attempts to address these problems have resulted in the use of additional external circuits to fill and blank the pulses."}
{"text": "The present invention relates generally to presence systems.\nWithin a presence system, when there is a state change associated with a source of a presence state, watchers of the source are notified of the state change. By way of example, if a source accepts a telephone call, the presence state of the source may switch from indicating that the source is available to indicating that the source is on a call, and watchers of source may be notified that the source is on a call. Alternatively, if a source begins to participate in a scheduled meeting, watchers of the source may be notified that a meeting state indicates that the source is involved in a meeting. That is, a presence state change notification may be sent to watchers to indicate that a presence state of a watched user has changed.\nCalendaring changes, e.g., changes in meeting states, for sources or users tend to occur at substantially the same time. Often, calendaring changes occur at the top of an hour. As such, a relatively high number of presence state change notifications may generally need to be sent at approximately the top of an hour. Sending a relatively high number of presence state change notifications at approximately the same time often creates a performance spike in the presence server that lasts until a queue of notifications is emptied. The performance spike may slow the sending of notifications, thereby causing watchers to receive the notifications a significant amount of time after state changes take place."}
{"text": "Mechanical couplings for joining pipe elements together end-to-end comprise interconnectable segments that are positionable circumferentially surrounding the end portions of pipe elements. The term “pipe element” is used herein to describe any pipe-like item or component having a pipe like form. Pipe elements include pipe stock, pipe fittings such as elbows, caps and tees as well as fluid control components such as valves, reducers, strainers, restrictors, pressure regulators and the like.\nEach mechanical coupling segment comprises a housing having arcuate surfaces which project radially inwardly from the housing and engage plain end pipe elements, shoulder end pipe elements, or circumferential grooves that extend around each of the pipe elements to be joined. Engagement between the arcuate surfaces and the pipe elements provides mechanical restraint to the joint and ensures that the pipe elements remain coupled even under high internal pressure and external forces. The housings define an annular channel that receives a gasket or seal, typically an elastomeric ring which engages the ends of each pipe element and cooperates with the segments to provide a fluid tight joint. The segments have connection members, typically in the form of lugs which project outwardly from the housings. The lugs are adapted to receive fasteners, such as nuts and bolts, which are adjustably tightenable to draw the segments toward one another.\nFor installation ready couplings of the type disclosed in U.S. Pat. No. 7,086,131 to Gibb et al., the coupling segments are preassembled at the factory, i.e., bolted together, but supported on the seal in spaced apart relation, one of the features which allow the pipe elements to be inserted into the coupling without first disassembling it. To facilitate positioning of the pipe elements within the coupling so that their arcuate surfaces align with circumferential grooves in the pipe elements the seal will have a stop comprising a circumferential ring. The stop ring is positioned concentric with the seal between the sealing surfaces which engage the pipe elements. The stop ring projects radially inwardly toward the center of the seal. Pipe elements are inserted into the seal between the segments until they contact the stop ring. The bolts holding the segments together are then tightened, drawing the segments toward one another to compress the seal against the pipe elements and engage the arcuate surfaces with the grooves to effect a fluid tight mechanical joint.\nTypically, the stop ring is made of the same flexible, resilient material as the seal. One disadvantage to such a construction is that a stop ring made of a soft, flexible material such as an elastomer can be damaged when pinched between the ends of the pipe elements, for example, during a bending test of the pipe joint or a major seismic event, where the pipe elements held by the coupling are substantially angularly displaced relative to the coupling. The pinched portion of the stop may become deformed, which is considered unacceptable for some applications. There is clearly a need for a stop ring which does not suffer the disadvantages described above to promote safety in critical fire protection systems while allowing for the economic advantages permitted by the use of installation ready couplings."}
{"text": "Vending machines have long been popular for dispensing food snacks because of the convenience. At first, vending machines were limited to dispensing packaged goods but have now been developed for dispensing goods that are prepared at the vending machine such as hot coffee, and chocolate. The trend continues to provide vending machines that mix and make the food product on demand rather than merely dispensing prepared and stored products. By mixing and making the product at the vending machine, a wider selection of flavors and choices are possible than from a vending machine that stores limited inventory of prepared products.\nWhen such products are made on demand, valves need to control the addition of various flavors to the base ingredient. For example, flavored syrup added to a base ice-cream is usually provided in liquid form. Displacement pumps and control valves need to control the addition of such different flavored syrups. Back flow prevention valves are also needed to prevent air from back flowing into the supply line to contaminate the supply. Previous valves were overly complex and often resulted in intermixing of different flavors. Known selector valves that rotate from one inlet to another and have a null or off position tend to get gummed up and stuck when viscous sugary fluids are used as the sugar crystallizes in the valve.\nWhat is needed is a dispensing valve that can control the flow of multiple supplies such as different flavors and be closeable and provide for a quick cleaning of the common downstream passages that receive the flow of a different supply with each dispensing. What is also needed is a valve capable of supplying a plurality of liquids with no rotary selector valve that provides for backflow prevention for preventing crystallization of any liquid syrups."}
{"text": "The present invention relates to systems and methods for linking orders in electronic trading systems. More particularly, the present invention relates to systems and methods which enable traders to link trading of goods, services, financial instruments, and commodities in electronic trading systems.\nIn recent years, electronic trading systems have gained wide spread acceptance for trading of a wide variety of goods, services, financial instruments, and commodities. For example, electronic trading systems have been created which facilitate the trading of financial instruments and commodities such as stocks, bonds, currency, futures, oil, gold, pork bellies, etc. As another example, online auctions on the Internet have become popular markets for the exchange of services and both new and used goods. In one embodiment of systems for electronic trading of financial instruments, for example, a first trader may submit a “bid” to buy a particular number of 30 Year U.S. Treasury bonds at a given price. In response to such a bid, a second trader may submit a “hit” in response to the bid in order to indicate a willingness to sell bonds to the first trader at the given price. Alternatively, the second trader may submit an “offer” to sell the particular number of the bonds at the given price, and then the first trader may submit a “take” or “lift” in response to the offer to indicate a willingness to buy bonds from the second trader at the given price. In such trading systems, the bid, the offer, the hit, and the take (or lift) are collectively know as “orders”. Thus, when a trader submits a bid, the trader is said to be submitting an order.\nModern day trading includes not only the buying and selling of a single type of item, but also more complex transactions involving exchanges of a combination of the same or different types of items. For example, in a typical spread transaction, one bond may be sold and another bond may be purchased as part of a single transaction. The trading of combinations of items in this way facilitates arbitrage, hedging, and speculation.\nHowever, because such combinations of items may have very complex relationships, there is a need to automate the trading of combinations of items. Thus, it is an object of the present invention to provide systems and methods for linking orders in electronic trading systems."}
{"text": "Frequency division multiplexing enables the concurrent communication of multiple signals over the same physical medium. In a frequency division multiplexed system, signals are frequency-converted to an assigned frequency band prior to being transmitted over the physical medium. To enable recovering the signals at the receiver, each of the different signals is assigned to a different frequency band or bands. The receiver then separates the received composite signal into the various frequency bands, and then processes the signal received in one or more of the assigned frequency bands to recover the information contained in that signal. Conventional circuitry utilized for separating the frequency bands, however, is costly.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "Mammalian cathepsins are cysteine-type proteases involved in key steps of biological and pathological events. Cathepsins are considered tractable drug targets as it is feasible to inhibit their enzymatic activity with small molecules and are therefore of interest to the pharmaceutical industry. Cathepsins are mainly located in the acidic compartments of the cells, like lysosomes and endosomes. In addition, cathepsins are secreted and work in the extracellular space, as well as in the cell cytoplasm and in the nucleus. In particular cathepsin L has a broad cellular distribution in all these compartments. By the use of alternative translation start sides downstream from the first AUG, alternative Cat L forms are generated devoid of the leader sequence. The truncated Cat L proteins are directed to the cytoplasm and the nucleus. Based on its cellular location, Cat L performs different cell biological activities.\nData from LDLrec (low density lipoprotein receptor) and Cat L deficient mice highlight the role of cathepsin L in atherosclerosis, as these mice show a reduced atherosclerotic phenotype (Kitamoto et al., Circulation 2007, 115:2065-75). Likewise, Cat L deficient mice have less severe lesions in the elastase induced model of abdominal aortic aneurism (Sun et al., Arterioscler Thromb Vasc Biol. 2011, 31:2500-8). Cat L contributes to vascular lesion formation by promoting inflammatory cell accumulation, angiogenesis, and protease expression. It is involved in matrix degradation, e.g. elastin and collagen, as a secreted protease, in autophagic cell death as cytoplasmic proteases (Mahmood et al., J. Biol. Chem. 2011, 286:28858-66) or by processing transcription factors like Cux-1 as a nuclear protease (Goulet et al., Mol. Cell. 2004, 14:207-19; Goulet et al., Biol. Chem. 2006, 387:1285-93). Human vascular disease samples from atherosclerotic vasculature or AAA (abdominal aortic aneurism) patients show strong upregulation of Cat L in diseases tissue (Liu et al., Atherosclerosis 2006, 184:302-11).\nCytoplasmic variants of Cat L seem to play a key role in proteinuric diseases. The podocyte is a key cell type maintaining the barrier function of the glomeruli in the kidney. Proinflammatory signals like LPS (lipopolysaccharide) induce Cat L expression. Cytoplasmic Cat L cleaves proteins that regulate the placticity of the cytoskeleton: dynamin and synaptopodin. Cat L deficient mice show reduced proteinurea in models of acute proteinurea (Reiser et al., J. Clin. Invest. 2010, 120:3421-31; Yaddanapudi et al., J. Clin. Invest. 2011, 121:3965-80).\nCat L deficient mice show a reduced metabolic phenotype when challenged towards different diabetic condition. Part of the mechanism is the cleavage of the insulin-receptor on skeletal muscle cells (Yang et al., Nat. Cell. Biol. 2007, 9:970-7), but matrix degradation as well as cleavage of the Cux-1 and its role in leptin signaling also contribute to the metabolic functions of Cat L (Stratigopoulos et al., J. Biol. Chem. 2011, 286:2155-70).\nCat L has also been shown to be upregulated in a variety of cancers ranging from breast, lung, gastric, colon to melanomas and gliomas. The cellular functions of Cat L in mediating apoptosis, lysosomal recyling, and cell invasion, make inhibition of Cat L in cancer an attractive target. The decrease of cell-cell adhesion by Cat L can partly be explained by cleavage of E-cadherin (Gocheva et al., Genes Dev. 2006, 20:543-56). The cleavage of extracellular matrix can also release growth factors from the matrix to interact with cell surface receptors."}
{"text": "Fluorescent linear lamps are widely used in office, retail and manufacturing, repair shop environments, and other settings, and typically fit into lighting fixtures with rectangular sheet metal boxes that have socket lamp holders at opposite ends thereof to retain and energize the fluorescent linear lamps. Fluorescent linear lamps have numerous shortcomings including generally poor light color quality, sometimes noisy operation, relatively high energy consumption, inclusion of toxic mercury (which makes disposal of fluorescent linear lamps problematic), and relatively short lifespans. In contrast, LED linear lamps can be designed to have any desired light color (measured in Kelvins), are quiet, are more energy efficient, do not include toxic mercury, and last a long time with operations of up to 50,000 hours (versus 10,000 for conventional fluorescent linear lamps.)\nWhen they were first introduced, LED linear lamps were much more costly than fluorescent linear lamps. With prices down significantly, it now makes sense to install LED linear lamps for new construction instead of fluorescent linear lamps. However, in cases of retrofit applications where fluorescent linear lamp light fixtures are already installed, updating fluorescent linear lamp fixtures to accommodate LED linear lamps is not always easy, convenient, or cost effective when considering union electrician labor rates. For example, in some situations the fluorescent linear lamp light fixture may have a light box that is slightly larger or smaller than is typical, the light fixture may be outfitted with old ballasts, or the fluorescent linear lamp light fixture may be designed to hold 2 or 4 parallel fluorescent linear lamps whereas the user wishes to use a different number of LED linear lamps therein. Also, many of the existing fixtures have old or damaged lamp sockets which need replacing when using traditional lamps with bi-pin ends to reduce risk of arcing or intermittent problems in the future. So, when converting to LED lamps changing the lamp sockets is costly in time and materials. As will be described further below, the invention allows installation of the LED lamp into the existing light fixture without the use of traditional lamp sockets. The existing sockets can be removed and the LED lamp installed easily and quickly.\nThere are currently available kits for converting fluorescent linear lamp light fixtures to work with LED linear lamp. For example, with the Everline Dimmable 21.6 W 4000K 2′×2′ LED Retrofit Kit, each LED linear bulb is incorporated into its own LED lensed modules/light bar. The LED lensed modules/light bars need to be screwed (with self-tapping screws) to the back wall of the light fixture and then wires therefrom will be connected to a light control module that will replace the fluorescent light ballast.\nIn the Litetronics® LED troffer retrofit kit, three LED linear lamps come preinstalled and spaced apart on a rack. The rack will be screwed to the back wall of the light fixture. However, this design is bulky to ship and can be relatively costly. Moreover, it does not allow customization by the user to change the number of the LED linear lamps.\nThere accordingly remains a need for adapters to allow custom or standard LED linear lamps to standard LED linear lamps in lighting fixtures originally outfitted with fluorescent linear lamps."}
{"text": "The present disclosure relates to a vehicle, and more particularly, to a personalized route planning system therefore.\nVehicles often include computer-implemented mapping systems. The mapping systems typically include route planning applications to provide users with directions between different locations. The route planning application includes representations of roads and intersections and one or more algorithms to output a suggested route of travel. These algorithms can output routes depending upon user-selected parameters. For instance, a route planning application can enable a user to select a time efficient route, or a distance efficient route.\nOver the last several years, users have grown to rely increasingly on route planning applications. Personalized tailoring of such routes, however, has been deficient."}
{"text": "Conventionally, there is known an alarm illumination device which gradually increases the brightness of a light source from a predetermined time before a preset wake-up time such that a sleeping person can be comfortably awakened at the wake-up time (see, e.g., Japanese Patent Application Publication No. H4-264289). Also, there is known an illumination device which gently changes illuminance from low illuminance to medium illuminance, and rapidly changes illuminance from medium illuminance to high illuminance by using light of at least three types of illuminance, i.e., low illuminance, medium illuminance and high illuminance (see, e.g., Japanese Patent Application Publication No. H7-318670). According to the illumination device, it is possible to cause the biological rhythm of the sleeping person to enter into an active phase while guiding the sleep state of the sleeping person from a deep state to a shallow state.\nHowever, in general, since many people get up after the consciousness is gradually awakened, it takes some time until they actually get up from waking up although there are individual differences. Accordingly, for example, as in the above-mentioned illumination device, when rapidly changing illuminance from medium illuminance to high illuminance, a person who has not been sufficiently awakened during the illuminance change from low illuminance to medium illuminance may feel uncomfortable because of the rapid illuminance change thereafter."}
{"text": "Heretofore, an eye cup is mounted on a window of the eye piece portion of a view finder optical system of a camera. This eye cup is used to cut light coming from around the window of the eye piece portion, thus a user can see an image at the screen clearly through the view finder optical system. And also, the light coming into the camera's main optical system is cut by the eye cup being in contact with the eye. Thus, the camera's main optical system is not disturbed by the light.\nHowever, some of the users wear glasses when viewing an object. The eye cup has a contact surface with a shape that is only made in accordance with the eye.\nThat is, the shape of the contact surface is not made in accordance with the shape of the lens of the glasses. Thus, a clearance exists between the lens and the contact surface. When a man puts on the glasses to take a photograph, light coming through the clearance is reflected by the face of the lens, then the light enters the camera's main optical system through the window of the eye piece portion."}
{"text": "In the utility industry, locks of the plunger operated type are often used to protect meters, meter boxes and valves that control the supply of electricity or gas. A suitable \"key\" or operating tool is provided to service personnel to enable them to open the locks as required. However these keys are often lost or stolen, and over a period of time many of these keys find their way into the hands of unauthorized personnel, who use them to remove and reverse meters, or to short across meter terminals, or to open valves that have been locked closed for nonpayment of utility bills.\nAlso, utilities that have adjoining territories prefer to have locks and keys that are not interchangeable, so that keys lost by service personnel of one utility cannot be used to open the locks of the adjoining utility.\nVarious efforts have been made to provide locks that can be opened only with a special key, and keys that cannot be modified to open other locks. For example, in my U.S. Pat. No. 4,252,006 there is illustrated a plunger lock and key in which the engaging means between the fingers of the key and the lock plunger have a configuration such that \"generations\" of locks and keys may be provided, so that a later \"generation\" key will open earlier \"generation\" locks, but will not open later generation locks. However in some cases the keys can be modified by unauthorized persons to open later generation locks.\nIn my copending application Ser. No. 281,701 filed 07/09/81 there is illustrated a lock and key combination which provides an improved system of locks and keys of various levels, which utilizes a key as disclosed and claimed in this application."}
{"text": "Mobile device networks such as wireless telephone networks are presently limited in the amount of data that is accessible by a mobile device user in a timely fashion. Wireless local area networks (WLANs) are increasingly being deployed in such public places as coffee shops, airports, hotels, and conference centers as a way to provide larger amounts of data to a mobile device user. WLAN access offers an opportunity for service providers to gain revenues from data services and for users to enjoy wireless high-speed data access in public spaces. Mobile network operators are interested in this opportunity because they already possess an established subscriber base with whom they presently have a billing relationship.\nBecause a public WLAN is not always operated by a mobile device user's own network, a protocol is required to authenticate a user across data networks. Authentication of a mobile device user is typically performed using Signaling System 7 (SS7) formatted communications between the mobile device network and the mobile device. However, communications between various networks takes place using the Internet Protocol (IP). SS7-format communications are not interchangeable with IP-format communications, making it difficult to implement a SS7-based authentication process using IP-format communications."}
{"text": "1. Field of the Invention\nThe present invention relates to bio-analysis, and more particularly a bio-analysis system integrating sample preparation process, and more particularly to a multi-channel bio-analysis system integrating sample preparation process.\n2. Description of Related Art\nBioanalysis, such as DNA analysis, is rapidly making the transition from a purely scientific quest for accuracy to a routine procedure with increased and proven dependability. Medical researchers, pharmacologists, and forensic investigators all use DNA analysis in the pursuit of their tasks. Yet due to the complexity of the equipment that detects and measures DNA samples and the difficulty in preparing the samples, the existing DNA analysis procedures are often time-consuming and expensive. It is therefore desirable to reduce the size, number of parts, and cost of equipment, to ease sample handling during the process, and in general, to have a simplified, low cost, high sensitivity detector.\nOne type of DNA analysis instrument separates DNA molecules by relying on electrophoresis. Electrophoresis techniques could be used to separate fragments of DNA for genotyping applications, including human identity testing, expression analysis, pathogen detection, mutation detection, and pharmacogenetics studies. The term electrophoresis refers to the movement of a charged molecule under the influence of an electric field. Electrophoresis can be used to separate molecules that have equivalent charge-to-mass ratios but different masses. DNA fragments are one example of such molecules.\nThere are a variety of commercially available instruments applying electrophoresis to analyze DNA samples. One such type is a capillary electrophoresis (CE) instrument. By applying electrophoresis in a fused silica capillary column carrying a buffer solution, the sample size requirement is significantly smaller and the speed of separation and resolution can be increased multiple times compared to the slab gel-electrophoresis method. These DNA fragments in CE are often detected by directing light through the capillary wall, at the components separating from the sample that has been tagged with a fluorescence material, and detecting the fluorescence emissions induced by the incident light. The intensities of the emission are representative of the concentration, amount and/or size of the components of the sample. In the past, Laser-induced fluorescence (LIF) detection methods had been developed for CE instruments. Fluorescence detection is often the detection method of choice in the fields of genomics and proteomics because of its outstanding sensitivity compared to other detection methods.\nHeretofore, CE instruments are designed to work with samples first prepared at other devices, and then loaded onto a sample tray in the CE instruments. Some of the sample preparation procedures could be quite involved, requiring manual and/or automatic procedures. Dedicated devices and systems have been designed to handle only sample preparation, involving steps such as sample extraction, purification, amplification, stabilization, etc., to produce samples that are suitable for separation by the CE instruments. For example, DNA samples may have to be prepared by a polymerase chain reaction (PCR) process, to amplify sufficient quantities of samples from a trace amount of DNA samples. The product of the PCR process may be subject to a CE process to verify the integrity or state of the result of the PCR process. The transfer from separately prepared samples to CE separation/analysis instruments require significant manual interventions, which affect overall throughput.\nIt would be desirable to develop a fully integrated bio-analysis system including built-in sample preparation process capabilities, to avoid user intervention during sample preparation and separation/analysis."}
{"text": "The present invention relates to the field of folding doors with flexible door leaves. More specifically, the invention relates to a door comprising a door leaf which is at least partly made of a flexible cloth material and which is movable between a closed position and an open, folded position, in which the door leaf is folded around a plurality of folding lines extended between opposite side edges of the door leaf, a plurality of guide members which are connected to the opposite side edges in a spaced-apart relationship along the same, and two side frames which extend adjacent to a respective side edge for guiding the guide members. Such a door is known from e.g. EP 0 113 634. The invention also relates to a method for assembling such a door.\nSince the 1970s there has been a great need to use rapidly moving doors in buildings for industrial use. This applies to openings indoors as well as in external walls, where the door provides shielding between different activities or prevents draughts/heat losses. Presently, rolling doors with flexible door leaves are used for this purpose, which doors are rolled up on an overhead drive shaft and which can be provided with transverse wind reinforcements on the door leaf to counteract wind load. For security reasons, rolling doors can be provided with a safety edge protection, a drop protection, etc.\nAlongside the development of rolling doors, there has been a development in foldable doors according to the introductory paragraph, in which the door leaf is instead folded as it is lifted during the opening process. These door leaves, too, are often provided with transverse wind reinforcements, comprising beams or sections which are suitably connected to the flexible door leaf. The wind reinforcements also contribute to the lateral stability of the door leaf.\nThe lifting arrangements of known folding doors vary from case to case, but usually the door leaf is lifted with the aid of at least one pair of belts/wires in the lowermost section, so that the transverse sections are gradually gathered in a bundle when the door is opened.\nEP 0 113 634 describes a folding door with transverse reinforcement sections. Every other section, beginning with the lowermost one, is extended into the side frames and supports guide rollers which are guided by the side frames in the depth direction, i.e. perpendicular to the door opening. The intermediate sections are shorter and have no guide rollers. Three lifting belts, which run vertically along the door leaf, are each connected to the bottom section. When the belts are rolled up on a transverse overhead shaft, they pull the bottom section upwards, which in turn successively pulls the other sections upwards so that the door leaf is folded in horizontal folds. Since every other section lacks guide rollers and consequently is not guided by the side frames, in the open position these non-guided sections will hang like a cradle by the intermediary of two superjacent guided sections, so that the door leaf is folded like a concertina. By virtue of the fact that the belts run on the exterior of the door leaf and on one and the same side thereof, all the non-guided sections are forced to fall out on the opposite side of the door leaf during the opening motion. Thus, in this known door, the lifting belts ensure that the non-guided sections fall out in one and the same direction.\nFR-A1-2,706,941 describes a folding door which, in conformity with the door in EP 0 113 634, has transverse reinforcement sections of which only every other section is guided by the side frames, and where the intermediate sections are non-guided in order to fall out sideways when the door is being closed. However, edge guide members are lacking, and the two side edges of the door leaf hang essentially completely unguided in the depth direction, received in the side frames. In this door, too, the lifting belts are used to ensure that the non-guided sections fall out on one and the same side of the door leaf. The lifting belts are located adjacent to the side frames.\nFR-A1-2,722,531 describes a door in which all the transverse reinforcement sections run in one and the same relatively wide guide track in the side frames and where the lifting belts are attached to the second lowest section and run through special belt loops in every other section. These loops result in the sections with loops gathering in a first bundle during lifting, while the sections without loops gather in a second bundle, hanging from the first bundle. The loops ensure that the sections without loops fall out on one and the same side of the door leaf in connection with lifting. Extra safety belts begin operating if the regular belts should break. All belts are located in the door opening between the side frames.\nSE 454,526 describes a technique for achieving forced folding of a door leaf, which is divided into horizontal, mutually foldable sections. In an embodiment shown in that document, the door leaf is designed in the form of a unitary, flexible piece of cloth, where every other section beginning with the lowermost is stiffened at its vertical side edges by means of rigid side borders. Every such rigid side border is provided with an upper and a lower guide pulley, which guide pulleys have a constant vertical relative position. These two pulleys run in an associated vertical guide track formed in the stationary side frame of the door. The guide track opening facing the door opening is provided with flanged edges for retaining the guide pulleys in the guide tracks. Thus, there is a plurality of guide tracks in each side frame. The number of guide tracks in each frame equals the number of sections provided with rigid side edges. Thus, only two guide pulleys run in each guide track, and, as a result of the stiffening, the stiffened sections are always vertically orientated in line with their associated guide track, and no folding takes place of these sections in connection with lifting. More specifically, the stiffened sections function as essentially completely rigid sections. In one example, the door leaf has three stiffened and three non-stiffened door leaf sections; and consequently three parallel guide tracks in each side frame.\nIn SE 454,526 mentioned above, two wires or the like are fastened to the lowermost, stiffened section for lifting and folding the door leaf. During lifting, the non-stiffened sections will be folded in between the stiffened sections, which assume a position beside each other like books on a shelf. When the door leaf has been lifted completely, a concertina-like bundle is obtained where the vertical, stiffened sections stand next to each other in a respective guide track and each intermediate, flexible section is extended obliquely downwards from the top of a stiffened section to the bottom of an adjacent stiffened section. In the lifted position, the whole bundle hangs from the section to which the wires are fastened.\nKnown folding doors of the type mentioned above exhibit various drawbacks depending upon the design chosen.\nIn the cases where the lifting belts and any associated loops are placed on the door leaf itself, there is a risk that individuals and vehicles will get caught in and lifted with the door leaf during opening. Moreover, such a placement is not aesthetically pleasing. Making holes for the lifting loops results in indication of fracture/weakening of the door leaf and additional manufacturing costs. In addition, centrally located lifting belts require a horizontal drive shaft or the like above the door.\nAnother drawback of the prior art doors is that the folding of the door leaf is effected in a non-reliable manner, or in a manner resulting in undesired wear of the door leaf. For example, the door leaf can be folded either inwards or outwards depending on the current pressure difference. This may, for example, result in the door leaf wearing against the upper edge of the door opening and/or the belts.\nAny pressure differences are absorbed by the transverse reinforcement sections, which, consequently, are squeezed against the side frames. In that way, in some known doors, the side edges of the door leaf are squeezed between the sections and the frames, resulting in the door leaf wearing out.\nMost known folding doors of the type described by way of introduction have a relatively wide side frame in the depth direction (i.e. transversely of the door opening) for receiving the side edge of the door leaf. Such a wide side frame is required to prevent the door leaf from jamming in the side frame during opening and closing. One drawback of having a wide side frame is that the door leaf can move in the depth direction in an undesired manner in connection with pressure differences, resulting in an undesired ability to move in the depth direction in the closed position, a poor aesthetic impression, and incomplete sealing. Moreover, a wide frame requires a large installation area and is expensive and heavy to make and assemble. A particular drawback of the door according to SE 454 526, wherein each stiffened section runs in its own guide track, is precisely that the side edges become very wide and costly as the height of the door and the number of sections increase, since a separate guide track is required for every other section of the door leaf.\nThese and other drawbacks of the prior art will appear clearly below in connection with the description of the invention.\nIn order to reduce the above-mentioned drawbacks of the prior art, according to the invention a door is provided of the type stated by way of introduction, i.e. a door comprising a door leaf which is at least partly made of a flexible cloth material and which is movable between a closed position and an open, folded position, in which open position the door leaf is folded about a plurality of folding lines extended between opposite side edges of the door leaf, a plurality of guide members which are connected to the opposite side edges in a spaced-apart relationship along the same; and two side frames which are extended adjacent to a respective side edge for guiding the guide members. The door according to the invention is characterised in that each side frame defines at least a first and a second guide groove, that said guide members comprise, at each side frame, a first set of guide members running in the first groove only of the side frame, and a second set of guide members running in the second groove only of the side frame, and that the first and the second guide members are connected to the door leaf in such a way that the side edges, in the folded position of the door leaf, run back and forth between the first and the second guide groove with said folding lines defined by the guide members.\nA xe2x80x9cflexible cloth materialxe2x80x9d could be any suitable kind of cloth, fabric or sheet of a flexible, foldable material, which can be coated or uncoated.\nWhen the door according to the invention is being opened or closed, the first guide members run in the first groove only and the second guide members run in the second groove only. In each groove, the associated guide members will be successively brought together during the opening motion. As a result, the mutual distance between the first guide members as well as the mutual distance between the second guide members will decrease when the door opens. Although, at present, it is probably preferable to have two guide grooves only in each side frame, it is within the scope of the invention to add one or more supplementary guide tracks, but in such variants it is still the case that the guide members in the first and the second guide groove are mutually brought together during the opening motion.\nThe expression xe2x80x9cguide groovexe2x80x9d can refer to a physical channel or the like, but it can also be interpreted as an abstract term and shall be considered to include all variants where the side edges are provided with special guide devices or means for defining two separate, predetermined movement paths or tracks for the guide members. Usually, the two guide grooves, which are defined by the side frames, are juxtaposed transversely of the door opening, but it is also possible that this distribution in the side frame itself is in a direction parallel to the door opening. In the latter case, there must be special connection members between the guide members and the edges of the door leaf, so that the attachment points in the edges of the door leaf run along two parallel lines or paths spaced from each other transversely of the door opening. In one embodiment, the first and the second guide groove can, for example, each be formed as a physical channel, whose side walls achieve the guiding of the guide members. These channels can be open towards the door opening but, with suitable connection members between the guide members and the door leaf edges, it is possible to turn the openings of the channels away from each other, so that one opening faces the front of the door and the other opening faces the rear of the door. As an alternative to physical channels, each guide groove can instead be defined by a rod or the like fixedly arranged in the side frame with which the guide members engage slidably in a suitable manner.\nUsually, the door according to the invention would be orientated with vertical side frames and a vertically guided door leaf. However, it is within the scope of the invention to place the door horizontally instead, but to facilitate the description and definition of the invention, terms such as xe2x80x9cliftingxe2x80x9d, xe2x80x9cvertical side framesxe2x80x9d, etc. are used throughout this specification. Accordingly, if the door is to be placed lying down, these orientation-determining expressions should be interpreted to include the horizontal case as well.\nIt should be noted that the above-mentioned xe2x80x9cplurality of guide membersxe2x80x9d can comprise xe2x80x9cfurther guide membersxe2x80x9d in addition to said first guide members and said second guide members, for example special guide members at the closing edge of the door leaf. Even if the first and the second guide members are normally located alternatingly in the first and the second guide groove, there may be portions of the door leaf where two adjacent guide members are located in the same guide groove.\nSeveral advantages are achieved by the invention by the provision of the double guide tracks in the side frames, as well as by the distribution of the guide members in the same:\n1. A first advantage of double guide tracks is that the folding of the door leaf becomes much more exact and controlled in comparison with how the folding takes place in the known doors. A controlled folding in the side frames in turn leads to generally safer functioning with a reduced risk of a breakdown, and to a considerable improvement in the appearance of the door leaf during operation. Moreover, no special pre-folding members or wear protection is necessary.\n2. As mentioned above, the side frames of the prior art doors must often be wide in the depth direction of the door opening in order to prevent the door leaf from jamming during lifting. The door according to the invention does not have that problem. Accordingly, a second advantage of double guide tracks is that the depth of the frame can be reduced considerably. The depth of the frame is mainly determined by the size of the guide elements in the depth direction of the door, but also by the amount of space required for the side edge itself of the door leaf.\n3. A third advantage of double guide tracks and of the side frames actively influencing the folding in the direction desired is that all lifting members, such as belts or wires, can be located protected within the side frames. Unlike in known doors, the lifting members need not be mounted on the surface of the door leaf for guiding the folding, but can be located protected in the side frames. This in turn means that both the lifting members and the environment are protected. The general appearance of the door also becomes more attractive with concealed lifting members. The driving can be achieved with two lifting points only, and if a variant with a transverse drive axle is used, it can be made with a less substantial dimension. Placing all the lifting members in the side frames also yields the advantage that no transverse drive shaft is needed above the door since belt drums can be attached directly to the side frames. However, it should be noted that, for example, in connection with very wide and/or heavy doors, it might be necessary to provide supplementary safety belts/lifting belts in the middle to prevent deflection. However, unlike in the prior art, it is not necessary to use such an additional belt for guiding the folding, but only for reducing the stress on any fall-out-preventing means in the side frame.\n4. A fourth advantage of double guide tracks is that, in its closed position, the door leaf can be positioned centrally in the depth direction between the guide tracks. This results in improved sealing and appearance, reduces wear and provides a more compact frame. In particular, the side edges of the door leaf can be guided in separate sections in the depth direction for obtaining an exceedingly compact door in the depth direction.\nNormally, the transverse folding lines, or extensions thereof, of the door leaf, will intersect the guide grooves. Accordingly, if the door leaf is provided with a plurality of transverse reinforcement members, each of which is extended between an associated pair of guide members, extensions of these reinforcement members can intersect the guide grooves for defining the folding lines of the door leaf. In order to obtain a straighter door leaf in the closed position, all the reinforcement members, or at least the majority of them, can lie alternatingly on the one and on the other side of the door leaf. In a special case, the two lowermost reinforcement members can be located in the same guide groove.\nWith respect to the space requirement at the upper part of the door, it will be appreciated that, in principle, the space required in the depth direction for the reinforcement members, when these are piled on top of each other according to the invention in two guide grooves, is only half as large as in the prior art where they are piled up in one and the same channel.\nIn a preferred embodiment of the invention, the first guide groove and the second guide groove in each side frame comprise a first physical guide channel and a second physical guide channel respectively, which are open in the direction of the door opening and have a width in the depth direction which is adapted to the dimensions of the guide members in the same direction. In this connection, the guide members can consist of non-rotatable sliding members or rollers. However, it is essential that no large play is required in the depth direction between the guide members and the side walls of the guide channels.\nEach guide channel can be provided with a fall-out-preventing means for retaining the guide members. The guide members can be designed themselves to prevent a fall if a lifting member breaks.\nAccording to a first embodiment, each side frame is provided with a U-section, whose bottom wall partly covers the two guide channels in order to form fall-out-preventing means, extended along the frame and open towards the door opening. In this connection, this U-section can have a double function since the side edge of the door leaf can be inserted in and seal against the U-section. One part of this U-section can be detachable for installation and maintenance. This embodiment yields the advantage that both the guide members and the side edges of the door leaf have a limited ability to move in the depth directionxe2x80x94i.e. they have good guiding in the depth direction and that the door leaf is centred in the depth direction relative to the guide tracks.\nAccording to a second embodiment, each side frame comprises a bottom wall, a first outer side wall and a first partition which both extend from the bottom wall for defining said first guide channel, and a second outer side wall and a second partition which both extend from the bottom wall for defining said second guide channel, wherein said first and second partitions define therebetween a space which receives the side edge of the door leaf. In this embodiment, the partitions serve two purposes: they define the guide channels and they receive therebetween the side edge of the door leaf in order to guide it along the side frame. In this embodiment, said partitions and said outer side walls can be provided with fall-out-preventing flanges adapted to retain the guide members in the guide channels.\nThe door leaf can be formed optionally as a continuous piece or divided into sections held together with e.g. transverse reinforcement sections. The door leaf can be formed entirely of a flexible cloth material, but the invention will also work if some door leaf sections are rigid. More specifically, the door leaf can be lifted in such a way that every other section is not folded, and these section can be made of a more rigid or a completely rigid material, while the other sections which are folded must be made of a flexible material.\nPreferably, there is at least a first flexible pulling member, such as a belt, a wire, a chain or the like, in each side frame for guiding the movement of the door leaf. If the guide grooves are physical channels, the pulling members can suitably be located in the same. In one embodiment, a direct lifting force is applied to only a single guide member in each side frame, called a driven guide member. If the door leaf runs alternatingly between the two guide tracks all the way down to its closing edge, the lifting can be effected in the lowermost guide member. However, in some cases, there may be a special safety arrangement with a bottom section having a reduced weight. In such a case, the lifting be effected in the second lowest guide member instead. If, however, there is only one pulling member in each side frame, these members can consist of a continuous pulling member. Moreover, there can be double pulling members or more in each side frame.\nIn principle, the lifting force applied to the driven guide members can be transmitted to superjacent guide members in two different ways. Either the design is such that the guide members strike against each other during the lifting, so that the lifting force is transmitted directly in the side frame. Alternatively, transverse reinforcement sections are used which are of such thickness that they will strike against each other before the guide members strike against each other. In this case, the lifting forces are instead transmitted by the intermediary of the reinforcement sections and, specifically, in this connection, guide members in the form of rotatable rollers can be used, which may be problematic if the guide members are to abut against each other.\nIn one embodiment of the invention there may also be a further pulling member in each side frame which applies a direct lifting force to a second driven guide member, the first and the second driven guide members running in different grooves. If, for example, the second driven guide member is located closer to the closing edge of the door leaf, its pulling member can be driven a somewhat longer distance than the first pulling member for achieving an xe2x80x9cextra liftxe2x80x9d of the second driven pulling member during the opening motion of the door. This can, for example, be achieved by the use of larger diameters in the winding drums for the second pulling members and/or greater thickness in the latter. Another possibility is to lift the last section more at the end of the lifting motion by virtue of only the lower part of the pulling member having a substantially greater thickness or to mount a member on the lower part of the pulling member which gives it an extra lifting motion at the end of the opening motion. The advantage of such an extra lift is that the vertical dimensions of the door leaf in the open position can be further reduced.\nFor easy transportation and installation of the door, each side frame can be divided into a shorter top part and a longer bottom part. The top parts are made with such a length that all guide members, which are connected to the side edges of the door leaf, can be received in the top parts simultaneously. In this way, the whole door leaf, all the reinforcement sections, all the lifting members, the upper part of the frames as well as the drive unit can be pre-assembled at the factory and be delivered to the installation site as a single unit. The top parts with the guide members inserted therein are mounted to the bottom parts only at the location where the door is to be installed. In assembling the frame parts, the guide grooves are likewise assembled, and the guide members and the pulling member can then be inserted into the side frames and the door can be used directly.\nThese and other embodiments and advantages of the invention will appear from the claims and from the following detailed description of preferred embodiments."}
{"text": "1. Technical Field\nThe present disclosure generally relates to testing devices and, particularly, to a probe for testing a printed circuit board.\n2. Description of Related Art\nDevelopments in electronic technology have brought increases in the number of components arranged on the printed circuit boards. A printed circuit board having components thereon requires quality testing of electronic properties before release. In a commonly used testing method, a plurality of welding portions is formed on the conductive wires of the printed circuit board. However, each welding portion is wider than the conductive wire, whereby signal transmission of the conductive wires and, correspondingly, the testing precision may be compromised.\nIn another testing method, a weld bead is applied on the conductive wire, and then contacted with a probe. However, the probe often has a flattened surface or a plurality of protrusions with tips. The contact area between the probe and the weld bead may be relatively small, because the weld bead is substantially spherical, which decreases the testing precision. Furthermore, the tips of the probe may cause damage or break the weld bead or the conductive wire easily, which also decreases the testing precision.\nTherefore, there is room for improvement within the art."}
{"text": "1. Field of the Invention\nThe present invention relates to a device and a method for controlling a substrate processing apparatus that performs a specified process on a substrate, and to a storage medium that stores the control program. More particularly, the present invention relates to a method for controlling the transfer of the substrate.\n2. Description of the Related Art\nThe recent substrate processing apparatuses installed in the semiconductor factory mostly include a transfer mechanism transferring a substrate and two or more processing chambers performing a specified process on the substrate. In the substrate processing apparatus including a plurality of processing chambers, it is important to transfer a large number of substrates to the processing chambers in a way to increase the throughput of the substrate process and increase the productivity of the products.\nSome of the conventional substrate processing apparatuses operate as follows. Different substrates are processed in different processing chambers at the same time. The processing is done after sequentially transferring the different substrates to the different processing chambers through different transfer paths for the different substrates (hereinafter this transfer method is also referred to as “OR transfer”). Alternatively, the transfer paths for the different substrates are controlled so that the substrates are sequentially processed via two or more processing chambers (see, for example, Japanese Patent Laid-Open Application No. 63-133532). The substrates may thus be processed efficiently.\nOthers of the conventional substrate processing apparatuses refer to a signal indicating which processing chamber is operatable, and transfer the substrate only to a group of the operatable processing chambers (see, for example, Japanese Patent Laid-Open Application No. 11-067869). If any of the processing chambers cannot operate due to a failure or the like, other processing chambers may be used to process the substrate efficiently."}
{"text": "Thanks to success of long term evolution (LTE)/LTE-advanced (LTE-A) for 4G mobile communication, interest in future mobile communication, that is, 5G mobile communication, is increasing and studies thereon are continuing.\nIn next-generation mobile communication, that is, 5G mobile communication, a data service having a minimum speed of 1 Gbps seems to be realized.\nIn 5G mobile communication, a turbo code, a polar code, a low density parity check (LDPC) code, etc. are considered as a channel coding method. Thereamong, the polar code is used by combining successive cancelation (SC) decoding and list decoding. However, in the list decoding, there is no method of excluding an erroneous decoding path through error correction before final decoding ends. Accordingly, the erroneous decoding path finally remains, thereby increasing an error probability and decreasing list gain."}
{"text": "With increasing needs to wirelessly identify and capture data related to physical objects, there has been an increased use of radio-frequency identification (RFID) and barcode technologies. RFID technologies offer the promise of non-line-of-sight identification, reading of many objects simultaneously with a single reader and the promise of supporting greater functionality including sensor inputs and large amounts of rewritable memory. For these reasons RFID has emerged as a preferred automatic identification and data capture technology.\nRFID readers or interrogators are used for communicating with and optionally powering tags affixed to objects. Readers can be found in fixed or mobile configurations and may be embedded in a variety of devices, including mobile phones. A unique ID code stored in the tag is communicated to the reader and associated with information in a database. At minimum an RFID tag consists of an electronic circuit attached to an antenna on some substrate. The electronic circuit may be chipless, containing only passive elements (e.g. inductors, capacitors, resistors, diodes), or may contain an integrated circuit containing active electronic devices such as transistors. A reader consists of an RF transceiver unit attached to one or more antennas.\nAt lower operating frequencies (such as 125 KHz or 13.56 MHz) where the wavelength of the electromagnetic field is large compared to the operating distance, the coupling between the reader and the tag is often described as “near-field” coupling. In this case, the reader antenna and the tag antennas are coupled together either inductively with coils or capacitively with planar electrodes, in which the tag antenna can be resonant or non-resonant.\nAt higher operating frequencies (such as 915 MHz or 2.45 GHz) where the wavelength of the electromagnetic field is small compared to the operating distance (generally 1 meter or more), the interaction between the reader and the tag is known as “far-field.” In this case, the tag communicates information to the reader by reflecting or scattering back some of the electromagnetic field that is incident on the tag. The amount of power that the tag is able to scatter back to the reader is dependent on the antenna shape, size, and tuning, and is generally described by a normalized parameter known as the scattering cross section.\nThe form of electromagnetic communication between a tag and a reader is important, since it influences the shape and form of the tag and reader antennas. Capacitively coupled antennas may be untuned or tuned and require two electrically disconnected electrodes (see e.g. U.S. Pat. No. 6,611,199). Inductively coupled antennas generally require tuned antennas in the form of a coil. Far-field antennas can employ either one or two electrodes, but also require tuning as well, in order to maximize the scattering cross-section.\nRFID tags of all forms present unique challenges to integration with products both electromagnetically and mechanically. Because RFID tags communicate with RFID readers via electromagnetic fields and waves, the product packaging materials and contents can strongly affect communications between reader and tag. Liquid and metallic materials are known to both absorb and reflect electromagnetic energy.\nThe physical integration of RFID tags with product packaging is generally challenging because many production processes have been well established and are unforgiving to significant changes or additions. Thus, there exists a need for innovative manufacturing methods as well as antenna designs to better integrate RFID tags into existing packaging materials and existing manufacturing methods at a low cost.\nIn the pharmaceutical industry, blister packages have emerged as a preferred method of packaging items for such reasons as security, product protection, display aesthetics, child-resistance, and medication compliance. Due to the increasing need to provide improved product tracking and tracing capabilities and additional security benefits there is a need to integrate RFID tags with blister packages.\nBlister packages have been in use since at least the 1960s for packaging of a variety of products, including items such as toys, tools, chewing gum, and medication. A large body of prior art exists which cover various blister package designs and packaging materials (see e.g. U.S. Pat. No. 3,054,502, U.S. Pat. No. 2,503,493, U.S. Pat. No. 3,380,578, and U.S. Pat. No. 5,954,204).\nDue to the existence of electrically conductive lidding film as a key component of the package, however, blister packages present unique challenges to integration of RFID tags. The lidding film typically consists of thin metal foil (0.6-1 mil=10 to 25 microns) and may incorporate additional laminated layers, such as paper or PET, and other coatings, for purposes such as heat sealing and printing. If RFID tag labels are applied to the lidding film, the electrical conductivity of the lidding film can detune RFID tag antennas and reflect electromagnetic fields and waves preventing necessary power transfer and communications—ultimately leading to poor RFID performance.\nOne solution to this problem is to add spacing or ferrite materials between the RFID tag and the blister package lidding film; however, this can add significant cost to overall package materials and production, while providing only a minor improvement in RFID performance.\nAdditionally, because the lidding film is designed to seal the contents in the package and act as a protective barrier, simply replacing the lidding film with any RFID tag label is not a viable option. The tag antenna must be designed so that the protective barrier is not compromised. Although various non-metallic non-conductive films have been developed as an alternative to metal foil lidding materials, the metal foils remain the most attractive in terms of cost.\nSince the blister seal must not be compromised and also since the surface area between adjacent blisters is limited, an additional challenge is to create an RFID antenna that will fit within the limited available area in between adjacent blisters.\nIn the case of the tuned coil antenna employed for inductive coupling, it is desirable to maximize the enclosed area of the coil as well as the number of turns in order to maximize the mutual inductance between the tag and reader and also achieve the proper inductance value to enable resonant tuning. In addition, since the thickness of the blister pack foil lidding is generally 30 microns or less, it is necessary to maintain the width of the metal coil traces to a few hundred microns or greater in order to prevent excessive resistive loss in the tag antenna coil. As a result of all these factors, it is a great challenge to create an antenna that will fit within the limited surface area of the blister pack and also avoid cutting the portions of the foil which seal the blisters.\nAlthough there have been several attempts to integrate blister packaging with RFID functionality, these prior inventions rely on antennas and electronic devices that are external to the blister package, and are not an integral part of the blister pack materials themselves. U.S. Pat. No. 6,244,462, for example, describes an external paper box or sleeve, with conductive traces, monitoring circuit, and transceiver into which a conventional blister package is inserted. Other prior art is specifically intended for monitoring the dispensing of medication in unit-dose blister packages, incorporate conductive traces located above the enclosed contents, which when broken, provide an indication that the contents have been removed. U.S. Pat. No. 6,574,166, for example, describes a package for monitoring medication compliance, where the conductive traces for sensing are integrated within the blister package, but the monitoring circuit and transceiver with antenna are located either external to the package or added as extra components to the package itself. The use of external devices or high-conductivity printed layers adds undesirable cost and complexity to the process of blister pack manufacture.\nIn order to achieve low-cost and large-scale manufacture of blister packages with RFID functionality, there still exists a need for innovative package designs and manufacturing methods which can enable better integration of an RFID circuit and antenna with the existing materials and processes used in blister pack manufacture."}
{"text": "Current methods of controlling transmitters, e.g. prototype, system under test, or experimental transmitters, in lab testing consist of an operator manually activating a control button, which provides a transmitter a required input to begin a transmission sequence. An operator would then observe a standard multi-meter output and record, by hand, a date/time a receive event occurred. This method of testing is both wasteful and inaccurate. A need presently exists for a way to automate and monitor a transmission sequence and receiver event. Another aspect is providing a monitoring and automation system which is flexible enough and capable of monitoring a variety of transmission sequence and/or receiver events to include particular types of timing or signal events."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical disk and a method for producing the same. In particular, the present invention relates to an optical disk in which a substrate on an incident side of laser light is thinner, and a method for producing the same.\n2. Description of the Related Art\nIn recent years, various studies have been conducted on the recording of optical information in the field of information recording. Recording of optical information can be conducted at higher density, and optical information can be recorded/reproduced in a non-contact manner; therefore, as a method for realizing the recording/reproducing of optical information at a low cost, applications for use in a wide range are being realized. Examples of current optical disks include those having a structure in which an information layer is provided on a transparent resin substrate with a thickness of 1.2 mm and protected by an overcoat or those having a structure in which an information layer is provided on one side or both sides of a transparent resin substrate with a thickness of 0.6 mm, and two substrates are attached to each other.\nRecently, in order to increase the recording density of an optical disk, a method for increasing a numerical aperture (NA) of an objective lens, a method for shortening a wavelength of the laser to be used, and the like have been considered. As the thickness of a recording/reproducing side substrate (i.e., substrate on an incident side of laser light) becomes smaller, the influence of aberration on a laser spot can be decreased, and an allowable value of a tilt of a disk can be increased. Because of this, it is proposed that the thickness of a recording/reproducing side substrate, a NA, and a laser wavelength are prescribed to be about 0.1 mm, about 0.85, and about 400 nm, respectively.\nIn a current DVD (digital versatile disk), mainly, a method is used in which two transparent resin substrates (thickness: 0.6 mm), on which film formation and the like are conducted, are attached, with radiation curable resin. Even when the thickness of a recording/reproducing side substrate becomes about 0.1 mm for the purpose of achieving high density, it is desirable to attach substrates to each other by the same method using the same facility as those currently used.\nHowever, with an optical disk in which two substrates are attached to each other, it is necessary to enhance durability. Furthermore, when the centers of two substrates are shifted from each other, deflections occur when the optical disk is rotated. Therefore, it is required to align the centers of two substrates with each other with high precision. There also is a demand for a method for easily producing such optical disks.\nTherefore, with the foregoing in mind, it is an object of the present invention to provide an optical disk that is recordable at high density by attaching two substrates to each other, and a method for producing the same.\nIn order to achieve the above-mentioned object, an optical disk of the present invention includes a first substrate having a signal area on a principal plane and a central hole A and a second substrate that is transparent and attached to the first substrate, wherein the second substrate is thinner than the first substrate and has a central hole B whose diameter is larger than that of the central hole A, and the first substrate and the second substrate are attached to each other with an adhesive member disposed therebetween so as to extend at least from an inner peripheral edge of the second substrate to an outer peripheral edge thereof.\nAccording to the above-mentioned configuration, an easy-to-handle optical disk is obtained that is capable of conducting high-density recording. Because of this, when a disk is handled, cracking or peeling of a contact portion can be prevented. The term xe2x80x9cradiationxe2x80x9d used herein includes a particle wave such as an electron beam and ultraviolet-rays and an electromagnetic wave.\nIn the above-mentioned optical disk, the adhesive member may be radiation curable resin. According to this configuration, an optical disk can be produced easily.\nIn the above-mentioned optical disk, a thickness of the second substrate may be in a range of 0.03 mm to 0.3 mm. According to this configuration, in particular, an optical disk that is recordable at high density can be obtained.\nIn the above-mentioned optical disk, the central hole B may be larger than a clamp area of the optical disk. According to this configuration, an optical disk can be fixed stably. Furthermore, when an optical disk is clamped, the second substrate can be prevented from peeling.\nIn the above-mentioned optical disk, the adhesive member may be disposed on an outer peripheral side of a clamp area or disposed so as to cover the entire clamp area. According to this configuration, since a thickness of the clamp area can be rendered uniform, a tilt is prevented from occurring during recording/reproduction.\nIn the above-mentioned optical disk, a thickness of a clamp area of the first substrate may be in a range of 1.1 mm to 1.3 mm.\nIn the above-mentioned optical disk, the first substrate includes, on the principal plane, at least one selected from the group consisting of a convex portion formed in a circular shape so as to surround the central hole A and having an outer diameter equal to or smaller than a diameter of the central hole B, and a concave portion formed in a circular shape so as to surround the central hole A and having a diameter equal to or smaller than the diameter of the central hole B.\nIn the above-mentioned optical disk, a height of the convex portion may be larger than a sum of a thickness of the second substrate and a thickness of the adhesive member.\nIn the above-mentioned optical disk, an average thickness of the adhesive member may be in a range of 0.5 xcexcm to 30 xcexcm.\nIn the above-mentioned optical disk, the optical disk is adapted for reproduction of information by application of a laser having a wavelength of 450 nm or less. According to this configuration, in particular, information can be recorded at high density.\nFurthermore, a first method for producing an optical disk of the present invention including a first substrate having a central hole A and a second substrate that is transparent and has a central hole B whose diameter is larger than that of the central hole A includes the processes of: (a) bringing the first substrate having a signal area on a principal plane and the second substrate that is thinner than the first substrate into contact with each other with radiation curable resin interposed therebetween so that the principal plane faces inside; and (b) irradiating the radiation curable resin with radiation to cure the radiation curable resin, thereby attaching the first substrate to the second substrate, wherein, in the process (a), the radiation curable resin is disposed so as to extend at least from an inner peripheral edge of the second substrate to an outer peripheral edge thereof.\nAccording to the first production method, an easy-to-handle optical disk that is recordable at high-density can be produced easily.\nIn the first production method, a thickness of the second substrate may be in a range of 0.03 mm to 0.3 mm.\nIn the first production method, the process (a) may include interposing the radiation curable resin between the first and second substrates, and rotating the first and second substrates to draw the radiation curable resin. According to this configuration, the thickness of resin easily can be rendered uniform.\nIn the first production method, the process (a) may include pouring the radiation curable resin onto the first substrate, rotating the first substrate to coat the first substrate with the radiation curable resin, and bringing the first substrate and the second substrate into contact with each other with the radiation curable resin interposed therebetween.\nIn the first production method, in the process (a), the first substrate and the second substrate are brought into contact with each other in a vacuum atmosphere. According to this configuration, air bubbles can be prevented from entering between the first substrate and the second substrate. The term xe2x80x9cvacuum atmospherexe2x80x9d as used here refers to an atmosphere with a reduced pressure (e.g., an atmosphere of 1000 Pa or less).\nIn the first production method, the first substrate may include, on the principal plane, at least one selected from the group consisting of a convex portion formed in a circular shape so as to surround the central hole A and having an outer diameter equal to or smaller than a diameter of the central hole B, and a concave portion formed in a circular shape so as to surround the central hole A and having a diameter equal to or smaller than that of the central hole B.\nIn the first production method, a height of the convex portion may be larger than a sum of a thickness of the second substrate and a thickness of the radiation curable resin.\nFurthermore, a second method for producing an optical disk of the present invention includes the processes of: (A) bringing a first substrate having a signal area on a principal plane and a central hole A and a second substrate that is transparent and thinner than the first substrate into contact with each other with radiation curable resin interposed therebetween so that the principal plane faces inside; (B) irradiating the radiation curable resin with radiation to cure the radiation curable resin, thereby attaching the first substrate to the second substrate; and (C) removing a part of the second substrate to form a central hole B whose diameter is larger than that of the central hole A in the second substrate, wherein, in the process (A), the radiation curable resin is disposed so as to extend at least from an outer periphery of a position where the central hole B is formed to an outer peripheral edge of the second substrate.\nAccording to the second production method, an easy-to-handle optical disk that is recordable at high density can be produced.\nIn the second production method, a thickness of the second substrate may be in a range of 0.03 mm to 0.3 mm.\nIn the second production method, the process (A) may include interposing the radiation curable resin between the first and second substrates, and rotating the first and second substrates to draw the radiation curable resin.\nIn the second production method, the process (A) may include pouring the radiation curable resin onto the first substrate, rotating the first substrate to coat the first substrate with the radiation curable resin, and bringing the first substrate and the second substrate into contact with each other with the radiation curable resin interposed therebetween.\nIn the second production method, in the process (A), the first substrate and the second substrate are brought into contact with each other in a vacuum atmosphere.\nFurthermore, a third method for producing an optical disk of the present invention includes the processes of: (i) opposing a first substrate in which a central hole A with a diameter dA is formed to a second substrate in which a central hole B with a diameter dB is formed with radiation curable resin interposed therebetween so that a center of the first substrate is aligned with a center of the second substrate; and (ii) irradiating the radiation curable resin with radiation to cure the radiation curable resin, wherein dA less than dB, and a thickness of the second substrate is in a range of 0.03 mm to 0.3 mm.\nAccording to the above-mentioned configuration, an optical disk that is recordable at high density can be produced with good precision.\nIn the third production method, in the process (i), the center of the first substrate is aligned with the center of the second substrate by using a pin that fits in the first and second central holes A and B. According to this configuration, it is easy to align the center of the first substrate with the center of the second substrate. As a result, an optical disk can be obtained in which deflections are unlikely to occur even when the optical disk is rotated at a high speed during recording/reproduction.\nIn the third production method, the process (i) may include the processes of: (i-1) fixing the second substrate on a table in which the pin is disposed so that the pin is inserted into the central hole B; (i-2) pouring the radiation curable resin onto the second substrate; (i-3) moving the first substrate so that the pin is inserted into the central hole A, thereby opposing the first substrate to the second substrate with the radiation curable resin interposed therebetween; and (i-4) rotating the first and second substrates to draw the radiation curable resin. According to this configuration, the thickness of the radiation curable resin can be rendered uniform. Therefore, an optical disk with good productivity and reliability can be produced.\nIn the third production method, the pin may include a first pin that fits in the central hole A and a second pin that fits in the central hole B, in the process (i-1), the second substrate may be fixed with the second pin, and in the process (i-3), the first substrate may be fixed with the first pin.\nThe third production method may include, after the process (i-1) and before the process (i-2), lowering an upper surface of the second pin below an upper surface of the second substrate.\nIn the third production method, the second pin may have a cylindrical shape, and the first pin may be inserted into the second pin.\nFurthermore, a fourth method for producing an optical disk of the present invention is a method for producing an optical disk including a first substrate in which a central hole A with a diameter dA is formed and a second substrate in which a central hole B with a diameter dB is formed, including the processes of: (I) coating at least one substrate selected from the group consisting of the first substrate and the second substrate with radiation curable resin; (II) opposing the first substrate to the second substrate with the radiation curable resin interposed therebetween in a vacuum atmosphere so that a center of the first substrate is aligned with a center of the second substrate; and (III) irradiating the radiation curable resin with radiation to cure the radiation curable resin, wherein dA less than dB, and a thickness of the second substrate is in a range of 0.03 mm to 0.3 mm.\nAccording to the fourth production method, an optical disk that is recordable at high density can be produced. Furthermore, the first substrate and the second substrate are opposed to each other in vacuum, so that air bubbles can be prevented from entering therebetween.\nIn the fourth production method, in the process (II), the center of the first substrate is aligned with the center of the second substrate by using a pin that fits in the first and second central holes A and B. According to this configuration, it is easy to align the center of the first substrate with the center of the second substrate.\nIn the fourth production method, the process (II) may include the processes of: (II-1) fixing the second substrate on a table in which the pin is disposed so that the pin is inserted into the central hole B; and (II-2) in a vacuum atmosphere, moving the first substrate so that the pin is inserted into the central hole A, thereby opposing the first substrate to the second substrate with the radiation curable resin interposed therebetween. According to this configuration, by fixing a second thin substrate on a table, the surface of the second substrate can be rendered flat; as a result, the thickness of the radiation curable resin can be rendered uniform. Furthermore, according to this configuration, air bubbles can be prevented from entering between the first substrate and the second substrate.\nIn the fourth production method, the pin may include a first pin that fits in the central hole A and a second pin that fits in the central hole B, in the process (II-1), the second substrate may be fixed with the second pin, and in the process (II-2), the first substrate may be fixed with the first pin.\nThe fourth production method further may include, after the process (II-1) and before the process (II-2), lowering an upper surface of the second pin below an upper surface of the second substrate.\nIn the fourth production method, the second pin may have a cylindrical shape, and the first pin may be inserted into the second pin.\nFurthermore, a production apparatus of the present invention is an apparatus for producing an optical disk including a first substrate in which a central hole A is formed and a second substrate in which a central hole B is formed, including: a coating member for coating at least one substrate selected from the group consisting of the first substrate and the second substrate with radiation curable resin; a disposing member for disposing the first substrate and the second substrate so that a center of the first substrate is aligned with a center of the second substrate; and an irradiating member for irradiating the radiation curable resin with radiation.\nAccording to the above-mentioned apparatus for producing an optical disk, the third and fourth production methods of the present invention can be conducted easily.\nIn the above-mentioned production apparatus, the disposing member may include a pin that fits in the first and second central holes A and B.\nIn the above-mentioned production apparatus, the pin may include a first pin that fits in the central hole A and a second pin that fits in the central hole B.\nIn the above-mentioned production apparatus, the second pin may have a cylindrical shape, and the first pin may be inserted into the second pin.\nIn the above-mentioned production apparatus, the disposing member may include a table for fixing the at least one substrate.\nIn the above-mentioned production apparatus, the disposing member further may include a container surrounding the table and an exhaust member for exhausting the container.\nThese and other advantages of the present invention will become apparent to those skilled in the art upon reading and understanding the following detailed description with reference to the accompanying figures."}
{"text": "The redox balance of cells is key to normal cell physiology. It is maintained by 3 systems: GSH/GSSG, NADPH/NADP; Thioredoxin (red)/Thioredoxin (oxd). Of these 3 systems, GSH/GSSG is the most widely studied for its implication in diseased states and for the development of rational therapeutic approaches (Townsend, A. J., Leone-Kabler, S., Haynes, R. L., Wu, Y., Szweda, L., and Bunting, K. D. (2001). Selective protection by stably transfected human ALDH3A1 (but not human ALDH1A1) against toxicity of aliphatic aldehydes in V79 cells. 130-132, 261-273). The diseased states associated with an imbalance in GSH/GSSG include major pathologies like cancers (Estrela, J. M., Ortega, A., and Obrador, E. (2006). Glutathione in cancer biology and therapy. Crit. Rev. Clin. Lab. Sci. 43, 143-181; O'Brien, M. L., and Tew, K. D. (1996). Glutathione and related enzymes in multidrug resistance. Eur. J. Cancer Oxf. Engl. 1990 32A, 967-978). Their one common aetiology is oxidative stress brought about by ROS and/or Reactive Nitrogen Species (RNS) that first cause a decrease in GSH due to the direct detoxification of ROS and RNS. This initial decrease in GSH is followed subsequently by a compensatory increase in GSH synthesis that cells bring into play in order to continue the detoxification of ROS/RNS and of newly formed electrophilic products such as 4-hydroxynonenal (HNE) and malondialdehyde (MDA)) produced by ROS attack on cellular lipids (Esterbauer, H., Schaur, R. J., and Zollner, H. (1991). Chemistry and biochemistry of 4-hydroxynonenal, malonaldehyde and related aldehydes. 11, 81-128).\nThe GSH paradox in cancer cells is that instead of the deficit in intracellular GSH that would have been expected, it is precisely the opposite that was found experimentally in many different cancer cells (Estrela, J. M., Ortega, A., and Obrador, E. (2006). Glutathione in cancer biology and therapy. Crit. Rev. Clin. Lab. Sci. 43, 143-181). But this increase in GSH has negative therapeutic repercussions as it protects cancer cells from chemo and radio therapies (Carretero, J., Obrador, E., Esteve, J. M., Ortega, A., Pellicer, J. A., Sempere, F. V., and Estrela, J. M. (2001). Tumoricidal activity of endothelial cells. Inhibition of endothelial nitric oxide production abrogates tumor cytotoxicity induced by hepatic sinusoidal endothelium in response to B16 melanoma adhesion in vitro. J. Biol. Chem. 276, 25775-25782).\nIn addition, if low levels of GSH must be obtained in cancer cells for chemotherapy to be effective, this is not the case for normal cells for not inducing collateral damage thereto.\nThe therapeutic approaches that are presently being used to lower cellular GSH in order to combat the chemoresistance of cancer cells, target GSH itself or the enzymes involved in GSH synthesis, GSH degradation and GSH efflux. There are already 10 GSH-lowering compounds that are in phases I, II and III of clinical trials as anticancer agents (Tew, K. and Townsend D (2011) Redox platforms in cancer drug discovery and development. Curr. Opin. Chem. Biol. 15, 156-161). They all have to be administered in combination with standard anti-cancer drugs, e.g. cyclophosphamide, taxol, vincristine, melphalan, etc.\nFurthermore, the enzymes targeted by these GSH-lowering drugs are those involved in GSH synthesis (gamma glutamyl cysteine ligase), GSH degradation (gamma-glutamyl transpeptidase) and GSH efflux (GSH-S-transferase). These same enzymes are however essential for protecting normal cells from ROS attack. Hence, there is a strong possibility of collateral damage to normal cells as the drugs cannot be delivered selectively to cancer cells and to cancer cells only.\nIn view of this, there is a need to find other therapeutic solutions which specifically and selectively target GSH in cancer cells.\nThe inventors of the present invention have unexpectedly found that a combination comprising an aminothiolester compound or a pharmaceutically acceptable salt thereof, in particular the S-methyl 4-(dimethylamino)-4-methylpent-2-ynethioate or a pharmaceutically acceptable salt thereof, and more particularly the 4-(Dimethylamino)-4-methyl-2-pentynethioic acid S-methyl ester fumarate, and a compound able to increase the H2O2 level in cancer cells of a subject, is useful as a medicament and able to treat cancer in a subject, wherein cancer cells of said subject do not overproduce H2O2. In particular, they found that a combination comprising an aminothiolester compound or a pharmaceutically acceptable salt thereof, in particular the S-methyl 4-(dimethylamino)-4-methylpent-2-ynethioate or a pharmaceutically acceptable salt thereof, and more particularly the 4-(Dimethylamino)-4-methyl-2-pentynethioic acid S-methyl ester fumarate, and a compound able to increase the H2O2 level in cancer cells of a subject, is useful as a medicament and able to treat cancer in a subject, wherein cancer cells of said subject do not overproduce H2O2 and have a level of GSH below 0.5 nmol for 25 000 cells.\nWithout being bound by any theory, when the compound able to increase the H2O2 level in cancer cells of a subject would have induced an increase in H2O2 level in the cancer cells, then the aminothiolester compound or a pharmaceutically acceptable salt thereof, would increase the levels of intra cellular metabolites that are produced by H2O2 attack, and, at the same time, GSH would thus be consumed in the detoxification of these electrophilic metabolites. As a result, insufficient GSH would be available in cancer cells to act as a scavenger of H2O2. Hence, levels of H2O2 should increase and should trigger-off the H2O2-dependent mechanisms in the mitochondrial (intrinsic) pathway of apoptosis.\nIn normal cells that have not undergone initially an attack by H2O2, intracellular GSH levels are already high (due to the absence of H2O2) so that the levels of any H2O2-induced electrophiles are below those in their cancer counterparts. However, H2O2 levels in normal cells can concomitantly rise when a compound able to increase the H2O2 level in cancer cells of a subject is used, if this compound is able to increase the level of H2O2 in both normal and cancer cells. In this latter case, H2O2 levels in normal cells will however still be lower than those in cancer cells and will thus remain below their apoptotic threshold upon treatment with the aminothiolester compound according to the invention or a pharmaceutically acceptable salt thereof."}
{"text": "This invention relates in general to actuating switches and in particular to a new and useful switch particularly for compressed gases in which during contact breaking the two switch parts are maintained in contact for a period of time and then the separation is effected with a driving separation which forces one of the parts to an end position.\nThe invention relates to an autopneumatic compressed gas switch having a piston-cylinder unit as a compression system for the quenching gas, a first contact piece and an insulating bushing surrounding the latter, which are firmly connected with the movable part of the piston-cylinder unit, and a second (counter-) contact piece which upon breaking contact is at first held in engagement with the first contact piece and follows it and then runs back under spring force.\nA similar compressed gas switch is known from German OS No. 29 18 145. In this compressed gas switch the second (counter-) contact piece is connected or latched directly with the first contact piece. A disadvantage in this solution is that the latch connection is arranged in a region in which the switching arc burns. It is therefore almost inevitable that the arc will impair the locking mechanism."}
{"text": "The present invention relates to a method for the closing, by laser, of packages of electronic circuits, notably hybrid circuit packages, said method minimizing mechanical stresses. It can be applied, more particularly, to the closing, by laser, of large-sized hybrid circuit packages, called macro-hybrid packages, where a metal lid is closed on a frame which is itself soldered to a ceramic substrate bearing screen-printed tracks and electronic components.\nMore generally, the invention can be applied to any hybrid circuit packages for which it is necessary to minimize the mechanical stresses related to closing by laser. These would be, for example, large-sized, all-metal solid packages.\nThe method of closing a metal package by laser is known to those skilled in the art. This method provides for tightly sealed soldering through the melting of the metal constituting the lid of the package. The parameters of the laser are determined as a function of the nature of the metals and the thickness of the lid. In the case of macro-hybrid packages for example, the substrate of which is constituted by a ceramic wafer supporting the hybrid electronic circuits, the frame is soldered to this ceramic substrate and not to a metal substrate. Experience has shown that any stress contributed by a solder made by means of a laser may prompt a fault at the solder joint between the frame and the ceramic. These stresses are actually related to the heating of the metals that get soldered to each other. The deformation thus created in the unit makes it impossible to obtain tight sealing in accordance with prevailing standards. Known methods consist, for example, in the soldering, transparently, of a lid made of an iron-nickel alloy to a frame made of an iron-nickel alloy or an iron cobalt-nickel alloy for example so that the laser beam, starting from a point of the perimeter of the macro-hybrid, travels along this perimeter and returns to the starting point. However, methods of this type lead to the above-mentioned deformations, prompted by the mechanical stresses exerted by the laser soldering. Indeed, before closing, the ceramic substrate of a macro-hybrid supporting the soldered frame generally displays a convexity of the order of 100 to 200 .mu.m due to the soldering with the frame whereas the macro-hybrid package, which has a surface area of the order of one dm.sup.2 once it is closed, displays a concavity of the order of 100 to 200 .mu.m corresponding to the stresses applied to the frame by the laser soldering. This results in defects of impervious sealing between the frame and the ceramic substrate, preventing any encapsulation in accordance with certain standards as laid down. Furthermore, when the laser soldering is done in the form of a continuous soldering, as described here above, the imperviousness between the lid and the frame is not acquired in a single pass: it is then necessary to double and even triple the seam of solder in the unsealed zones. This contributes greatly to increasing the stresses exerted by the soldering and increases manufacturing costs."}
{"text": "1. Field of the Invention\nThe present invention relates to photoelectric conversion apparatus and information processing apparatus, and more particularly to photoelectric conversion apparatus suitably applicable to input portions in facsimile devices, image readers, scanners, copiers, electronic blackboards, etc., which are information processing apparatus for reading image information while relatively moving the original or the like subjected to image reading in close contact on a primary line sensor, and the information processing apparatus therewith.\n2. Related Background Art\nRecently, in order to further decrease the size or further enhance the performance as to information processing apparatuses such as facsimile devices, image readers, etc., elongate line sensors that can be used in a 1:1 optical system have been developed as photoelectric conversion apparatuses for information processing apparatuses. In order to further decrease the size and the cost, photoelectric conversion apparatus and information processing apparatus have been developed for directly detecting reflected light from the original by a sensor without using a 1:1 fiber lens array but through a transparent spacer of glass or the like.\nFIG. 1 is a diagrammatic, perspective view of a conventional photoelectric conversion apparatus, and FIG. 2 is a diagrammatic cross section thereof when the photoelectric conversion apparatus is cut along 2--2 in FIG. 1. In FIG. 1 and FIG. 2, reference numeral 1 designates a photosensor, 101 a sensor substrate, 102 photoelectric conversion elements, 103 a wiring part, 104 a mount plate, and 105 a transparent protection layer.\nAs shown in FIG. 1 and FIG. 2, a plurality of photoelectric conversion elements 102 are arranged in line on the sensor substrate 101, and the transparent protection layer 105 comprised of a thin glass plate or the like is provided on the photoelectric conversion elements 102 in order to protect the photoelectric conversion elements 102 and to act as a spacer between the photoelectric conversion elements 102 and the original Output signals from the photoelectric conversion elements 102 are arranged as capable of being output to the outside through the wiring part 103. These sensor substrate 101, photoelectric conversion elements 102, wiring part 103, and transparent protection layer 105 are incorporated with the mount plate 104, thus composing the photoelectric conversion apparatus 1. The mount plate 104 is made of an easily-moldable material such as a resin, and is molded except for an optically transparent portion of the member.\nThe structure of a photoelectric conversion element used in such a photoelectric conversion apparatus is next explained briefly referring to FIG. 3.\nFIG. 3 is a diagrammatic cross section for illustrating an example of the photoelectric conversion element. In FIG. 3, reference numeral 170 is a substrate, 171 a light-shielding layer, for example, of aluminum or chromium, 172 an insulator layer, for example, of silicon nitride, 173 a semiconductor layer, for example, of i-type amorphous silicon, 174 an ohmic contact layer, for example, of n.sup.+ -type non-single-crystal silicon, 175 an electrode layer, for example, of aluminum, 176 a passivation layer, for example, of silicon nitride or polyimide, and 177 an adhesive layer, for example, of an epoxy resin.\nAs shown in FIG. 3, the sensor substrate 101 has a photoelectric conversion element having the light-shielding layer 171 provided on the substrate 170 in correspondence to a photoelectric conversion portion so as to prevent illumination light from the substrate side from entering the photoelectric conversion portion, the insulator layer 172 provided on the light-shielding layer 171, the semiconductor layer 173 provided on the insulator layer, and the electrode layer 175 provided through the ohmic contact layer 174 above the semiconductor layer with a space for an incident area of light, and the passivation layer 176 provided on the photoelectric conversion element; and the transparent protection layer 105 is provided through the adhesive layer 177 above the sensor substrate 101.\nLight L emitted from a light source (not shown) passes through the sensor substrate 101, adhesive layer 177, and transparent protection layer 105 to reach the original P, and reflected light from the original P is incident to the space formed in the electrode layer 175 to be photoelectrically converted according to the incidence of light.\nFIG. 4 is a diagrammatic, sectional, structural drawing for illustrating an example of a facsimile device with a communication function as an example of the information processing apparatus having the above photoelectric conversion apparatus.\nIn FIG. 4, reference numeral 180 denotes the facsimile device, 181 a feed roller, 182 a separator pawl, 183 a conveying roller, 184 a system control substrate, 185 a platen roller, 186 a recording head, 187 a power supply section, 188 the photoelectric conversion apparatus, 189 an operation panel, P the original, and W a recording medium.\nAs shown in FIG. 4, the facsimile device 180 supplies the original P to an image reading section by the feed roller 181 when the original P is inserted thereinto. In the original reading section the photoelectric conversion apparatus 188 is disposed and the original P is conveyed as urged against the reading portion by the conveying roller 183 opposed to the photoelectric conversion apparatus 188. This urging is effected by an urging device such as a spring, not shown. As urging the conveying roller 183 and/or the photoelectric conversion apparatus 188 against each other, the original P conveyed to between the conveying roller 183 and the photoelectric conversion apparatus 188 is further conveyed in a discharge direction. The separator pawl 182 is used for separating and feeding the originals P set in a pile one by one.\nFurther, the facsimile device 180 has the recording head 186 for recording an image received or information read by the above photoelectric conversion apparatus 188 in a recording medium W, and the platen roller 185 for conveying the recording medium W for recording information by the recording head 186.\nThe power supply 187 is a power-supply portion for driving the facsimile device 180, the system controller 184 is provided for controlling an image reading means including the photoelectric conversion apparatus 188 and the recording means including the recording head 186, and the operation panel 189 is a so-called control portion of the facsimile device 180.\nThe photoelectric conversion apparatus 188 is mounted, as shown, to a frame 4 (a main body frame in the drawing) provided in the facsimile device 180. This frame is normally provided in order to attain mainly the strength of the apparatus main body.\nFIG. 5 is a diagrammatic, structural drawing for illustrating a mounting portion of the photoelectric conversion apparatus 188 and surroundings thereof.\nIn FIG. 5, reference numeral 191 represents a light source such as light-emitting diodes, 192 a sensor frame, 193 a light-source substrate, and 194 a connector. For using the light-emitting diodes as a light source 191, a plurality of light-emitting diodes are arranged at intervals on the light-source substrate 193. The photosensor 1 and light-source substrate 193 are mounted to the sensor frame 192 so as to be incorporated therewith, thereby composing the photoelectric conversion apparatus 188. The connector 194 is provided for supplying the power for driving the photoelectric conversion elements or the light-emitting diodes and/or for outputting electric signals carrying information output from the photoelectric conversion elements.\nThe photoelectric conversion apparatus 188 is mounted by unrepresented means in a recess formed in the frame 4.\nFIG. 6 is a diagrammatic, perspective assembly drawing for illustrating the mounting portion of the photoelectric conversion apparatus 188 and surroundings thereof\nIn FIG. 6, reference numeral 195 stands for a flexible board for outputting signals from the photosensor 1, 196 a flexible board for outputting the signals from the photoelectric conversion apparatus 188 to the outside, and 197 an electric connection portion in FIG. 7 to be connected with the flexible board.\nAs shown in FIG. 6, the photoelectric conversion apparatus 188 is arranged so as to be dropped into the recess formed in the frame 4, and the flexible board 196 is connected to the electric connection portion on the apparatus body side, thereby electrically being connected with the apparatus body. The conveying roller 183 is disposed above the reading portion of the photoelectric conversion apparatus 188, as described above. The signal line from the photosensor 1 is first electrically connected through the flexible board 195 to the light-source substrate 193 and then electrically connected through a signal processing circuit provided in the light-source substrate 193 and through the flexible board 196 to a processing circuit on the apparatus body side.\nFIG. 7 is a diagrammatic, structural drawing for illustrating an example of the connection relation between the photoelectric conversion apparatus 188 and the system control substrate 184 in the facsimile device 180.\nAs shown in FIG. 7, the flexible board 196 from the photoelectric conversion apparatus 188 is electrically connected with the connector 197 formed in the system control substrate 184.\nHowever, with attempt to achieve further reductions of cost, size, and weight for the information processing apparatus constructed in the above structure, there was virtually no room for further reductions of cost, size, and weight, because the photoelectric conversion apparatus is constructed in a unit structure using the sensor frame. Another problem was that the limits of further reductions of size and weight would impose restrictions on freedom of engineering design of the apparatus main body or freedom of design thereof."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a method for directly writing data into an optic disk without a computer system; and in particular to a method for directly writing data retrieved from an electronic data storage into an optic disk without a computer system interfacing therebetween.\n2. The Related Art\nWith the rapid development and prevalence of electronic storages, such as compact flash memory, more and more data are stored in the electronic storage for portability and fast access. To more space-efficiently store the data, some people prefer to transfer the data from the electronic storage to optic disks for data backup purposes. Heretofore, the optic disk drive must be connected to a computer system for data writing operation. Thus, the data have to be read into the computer system and then written by the computer system into the optic disk accessed by means of the optic disk drive. This causes problems. For example, a computer system is a must in transferring data from a portable compact flash memory device to an optic disk. Thus, such a data transfer operation cannot be carried out without a computer system having proper data ports.\nThus, the present invention is aimed to solve the above problem by providing a method for directly transferring data from a compact flash memory device to an optic disk without a computer system interfacing therebetween.\nAccordingly, an object of the present invention is to provide a method for writing data retrieved from an electronic storage to an optic disk without a computer system interfacing therebetween.\nAnother object of the present invention is to provide an optic disk drive capable to perform a direct writing operation to an optic disk without being controlled by a computer system.\nTo achieve the above objects, in accordance with the present invention, there is provided a method for directly writing data into an optic disk that is performed by an optic disk drive incorporating a control unit to which an external data storage device, such as compact flash memory device, is connected. The method comprises steps of (1) initiating a writing operation, (2) setting the optic disk drive to busy condition, (3) checking if an optic disk is properly loaded and if the external memory device is correctly connected, (4) checking if the optic disk is a UDF disk; (5) issuing a warning, if it is not, (6) creating a folder in the optic disk, (7) retrieving data from the external data storage device and writing the data into the folder of the optic disk, and (8) ending the writing operation. No computer-based interface is required between the optic disk drive and the external data storage device in performing the data writing operation."}
{"text": "1. Field of Use\nThe present invention relates to cache systems and, in particular, to cache systems includable within minicomputer and microprocessing systems.\n2. Prior Art\nIt is well known in the art to provide cache systems within computer systems to improve overall system performance and provide for reliable operation. Examples of such systems are disclosed in U.S. Pat. No. 3,820,078 to John L. Curley, et al. and in IBM Technical Disclosure Bulletin titled \"Removal of Failing Buffer Sections in a Buffer Backing Store\" by M. W. Bee, et al., Vol. 13, No. 2, dated July 1970. In those systems, reliable operation is achieved by invalidating cache memory locations detected as having bad parity. This requires additional bits to be associated with the cache locations and can add considerably to the cost and complexity of the cache system. More importantly, it requires processing time for carrying out such invalidating operations.\nOther prior art cache systems permit the cache to be bypassed upon the detection of fault conditions by the central processing unit (CPU) associated therewith. An example of this type of system is disclosed in U.S. Pat. No. 4,195,343 to Thomas F. Joyce which is assigned to the same assignee as named herein. In general, this type of cache system is designed to report two types of errors to the CPU, a memory \"red\" error condition indicative of an uncorrectable error and a memory \"yellow\" condition indicative of a correctable error. Upon the receipt of a \"red\" error signal or the detection of a byte data parity error in received memory data, the CPU switches the entire cache off-line, reports the error to the operating system and continues processing.\nWhile the above systems allow disconnection of a cache unit as a consequence of a fault condition, it requires the CPU to process such fault conditions. This can prove time consuming and could also result in loss of valuable information since such diagnosis relies in part upon the types of error conditions reported by the cache system itself.\nAccordingly, it is a primary object of the present invention to provide a reliable cache system which is low in cost and has minimal complexity.\nIt is still a further object of the present invention to provide a cache system which has improved maintainability thereby increasing system reliability."}
{"text": "The present invention relates to surgical systems and, in various arrangements, to grasping instruments that are designed to grasp the tissue of a patient, dissecting instruments configured to manipulate the tissue of a patient, clip appliers configured to clip the tissue of a patient, and suturing instruments configured to suture the tissue of a patient, among others.\nCorresponding reference characters indicate corresponding parts throughout the several views. The exemplifications set out herein illustrate various embodiments of the invention, in one form, and such exemplifications are not to be construed as limiting the scope of the invention in any manner."}
{"text": "Solar energy generation is a rapidly growing technology worldwide and offers the potential of almost unlimited clean and sustainable energy. However, the use of solar electric technology has been limited by the costs associated with installing solar panels to existing and new structures and facilities.\nWhen installing a solar module on a rail, various clamps must be utilized due to the varying sizes of the modules and various rail configurations. As a result, it is desirable to have a clamp that can be used to secure different types of modules to different types of rails.\nThe solar module is often installed on a roof or other surface for exposure to sunlight. As a result, the installed solar module can be viewed. Accordingly, it is desirable to have a clamp that is aesthetically pleasing and is preferably hidden under the module frames.\nBecause maintenance may be required for the solar modules and because the solar modules may be installed on a roof or other surface where access is often needed, safety is also an important consideration. Modules can be installed at a variety of heights, commonly about three to eight feet off the ground, and on a variety of surfaces, such as a roof of a building. When someone is walking next to a solar module, it is desirable that the rail does not extend past the edge of the module."}
{"text": "1) Field of the Invention\nThe present invention relates to an improved microwave-interactive cooking package. In particular, the present invention relates to high efficiency, safe and abuse-tolerant susceptor and foil materials for packaging and cooking microwavable food.\n2) Description of the Related Art\nAlthough microwave ovens have become extremely popular, they are still seen as having less than ideal cooking characteristics. For example, food cooked in a microwave oven generally does not exhibit the texture, browning, or crispness that are acquired when food is cooked in a conventional oven.\nA good deal of work has been done in creating materials or utensils that permit food to be cooked in a microwave oven to obtain cooking results similar to that of conventional ovens. The most popular device being used at present is a plain, susceptor material, which is an extremely thin (generally 60 to 100 xc3x85) metallized film that heats under the influence of a microwave field. Various plain susceptors (typically aluminum, but many variants exist) and various patterned susceptors (including square matrix, xe2x80x9cshower flower,xe2x80x9d hexagonal, slot matrix and xe2x80x9cfusexe2x80x9d structures) are generally safe for microwave cooking. However, susceptors do not have a strong ability to modify a non-uniform microwave heating pattern in food through shielding and redistributing microwave power. The quasi-continuous electrical nature of these materials prevents large induced currents (so limiting their power reflection capabilities) or high electromagnetic (E-field) strengths along their boundaries or edges. Therefore their ability to obtain uniform cooking results in a microwave oven is quite limited.\nElectrically xe2x80x9cthickxe2x80x9d metallic materials (e.g., foil materials) have also been used for enhancing the shielding and heating of food cooked in a microwave oven. Foil materials are much thicker layers of metal than the thin, metallized films of susceptors. Foil materials, also often aluminum, are quite effective in the prevention of local overheating or hot spots in food cooked in a microwave by redistributing the heating effect and creating surface browning and crisping in the food cooked with microwave energy. However, many designs fail to meet the normal consumer safety requirements by either causing fires, or creating arcing as a result of improper design or misuse of the material.\nThe reason for such safety problems is that any bulk metallic substance can carry very high induced electric currents in opposition to an applied high electromagnetic field under microwave oven cooking. This results in the potential for very high induced electromagnetic field strengths across any current discontinuity (e.g., across open circuit joints or between the package and the wall of the oven). The larger the size of the bulk metallic materials used in the package, the higher the potential induced current and induced voltage generated along the periphery of the metallic substance metal. The applied E-field strength in a domestic microwave oven might be as high as 15 kV/m under no load or light load operation. The threat of voltage breakdown in the substrates of food packages as well as the threat of overheating due to localized high current density may cause various safety failures. These concerns limit the commercialization of bulk foil materials in food packaging.\nCommonly owned Canadian Patent No. 2196154 offers a means of avoiding abuse risks with aluminum foil patterns. The structure disclosed addresses the problems associated with bulk foil materials by reducing the physical size of each metallic element in the material. Neither voltage breakdown, nor current overheat will occur with this structure in most microwave ovens, even under abuse cooking conditions. Abuse cooking conditions can include any use of a material contrary to its intended purpose including cooking with cut or folded material, or cooking without the intended food load on the material. In addition, the heating effectiveness of these metallic materials is maximized through dielectric loading of the gaps between each small element that causes the foil pattern to act as a resonant loop (albeit at a much lower Q-factor (quality factor) than the solid loop). These foil patterns were effective for surface heating. However, it was not recognized that a properly designed metallic strip pattern could also act to effectively shield microwave energy to further promote uniform cooking.\nCommonly owned U.S. Pat. No. 6,133,560 approaches the problem differently by creating low Q-factor resonant circuits by patterning a susceptor substrate. The low Q-factor operation described in U.S. Pat. No. 6,133,560 provides only a limited degree of power balancing.\nThe present invention relates to an abuse-tolerant microwave packaging material which both shields food from microwave energy to control the occurrence of localized overheating in food cooked in a microwave, and focuses microwave energy to an adjacent food surface.\nAbuse-tolerant packaging according to the present invention includes one or more sets of continuously repeated microwave energy interactive/reflective segments disposed on a microwave-safe substrate. Each set of reflective segments defines a perimeter equal to a predetermined fraction of the effective wavelength in an operating microwave oven. Methodologies for choosing such predetermined fractional wavelengths are discussed in U.S. Pat. No. 5,910,268, which is incorporated herein by reference. The reflective segments can be metallic foil segments, or may be segments of a high optical density evaporated material deposited on the substrate. The terms xe2x80x9cfractionxe2x80x9d or xe2x80x9cfractionalxe2x80x9d as used herein are meant in their broadest sense as the numerical representation of the quotient of two numbers, i.e., the terms include values of greater than, equal to, and less than one (1).\nIn a first embodiment, the length of the perimeter defined by a first set of microwave energy interactive/reflective segments is preferably approximately equal to an integer multiple of the effective wavelength of microwaves in an operating microwave oven, such that the length of the perimeter is resonant with the effective wavelength. In a second embodiment, the length of the perimeter defined by the reflective segments is approximately equal to an integer multiple of one-half the effective wavelength of microwaves in an operating microwave oven, such that the length of the second perimeter is quasi-resonant with the effective wavelength.\nEach segment in the first set is spaced from adjacent segments so as to create a (DC) electrical discontinuity between the segments. Preferably, each first set of reflective segments defines a five-lobed flower shape. The five-lobed flower shape promotes uniform distribution of microwave energy to adjacent food by distributing energy from its perimeter to its center.\nPreferably, abuse-tolerant packaging according to the present invention includes a repeated second set of spaced microwave energy interactive/reflective segments that enclose each first set of reflective segments and define a second perimeter. In the first embodiment, this second perimeter preferably has a length approximately equal to an integer multiple of the effective wavelength of microwaves in an operating microwave oven, such that the length of the second perimeter is resonant with the effective wavelength. In the second embodiment, this second perimeter preferably has a length approximately equal to an integer multiple of one-half the effective wavelength of microwaves in an operating microwave oven, such that the length of the second perimeter is quasi-resonant with the effective wavelength.\nA third embodiment of abuse-tolerant packaging according to the present invention includes, in addition to the second set of reflective segments, a repeated third set of spaced microwave energy interactive/reflective segments that enclose each second set of reflective segments and define a perimeter approximately equal to another predetermined fraction of the effective wavelength of microwaves in an operating microwave oven.\nFurther embodiments of the invention may be created by varying the shapes of the perimeters formed by the reflective segments, while maintaining the desired predetermined fraction of the effective wavelength for the length of the perimeters. Appropriate shapes within the scope of the present invention may be, for example, circles, ovals, and other curvilinear shapes, triangles, squares, rectangles, and other polygonal shapes. Curvilinear shapes are preferably symmetrical to aid in the assembly of shapes in an array. Similarly, polygonal shapes are preferably right and equilateral polygons to help in the formation of nested arrays of the shapes."}
{"text": "Field\nEmbodiments described herein relate to an electrically-rewritable nonvolatile semiconductor memory device.\nDescription of the Related Art\nA NAND type flash memory has a memory cell array which is configured by arranging memory strings each including a plurality of memory cells connected in series. Both ends of each memory string are connected to a bit line and a source line through select transistors respectively. The control gate electrodes of the memory cells of each memory string are connected to different word lines respectively. In each memory string, the plurality of memory cells are connected in series with sources and drains shared between them. The NAND type flash memory can have a small unit memory cell size, because select transistors and their bit line contacts and source line contacts are shared among a plurality of memory cells. Further, the NAND type flash memory is suitable for miniaturization, because the word lines and the device regions of the memory cells have a shape that resembles a simple stripe shape, and hence a flash memory having a large capacity is realized.\nAs the miniaturization of NAND type flash memories progresses, interference between adjoining cells and influence due to elapse of time after data writing increase, which might change the memory cell data. For example, when data written in a memory cell remains un-accessed for a long time, there occurs a phenomenon that electrons are discharged from the charge accumulation layer of the memory cell, changing the threshold voltage of the memory cell to a lower value. Hereinafter, this phenomenon will be referred to as deterioration of data retention. If deterioration of data retention occurs, a data reading operation may fail."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nCasting processes for forming articles using molds and cores employ casting chambers including outer molds and inner core elements each having features and reliefs that form details, recesses, and cavities in a casting when molten material such as liquid metal is poured into the mold. One casting formed by such a casting process is an engine block formed from molten cast iron or molten aluminum alloys. Inner core elements can be constructed from bonded sand. The inner core elements are extracted from the casting subsequent to the forming process. Portions of the casting may be subject to high-stress in-use, and it may be desirable to impart varying metallurgical properties to those portions. For example, a time-rate of removal of thermal energy from liquid metal during casting affects grain structure, with increased cooling and solidification of the poured liquid metal leading to an improvement, in general, of material properties such as tensile strength, fatigue strength, and in some cases machinability.\nKnown casting processes use thermal chill devices in proximity to specific portions of a casting in place of or in conjunction with features on the mold and core elements. This includes using chill devices at bulkheads and crankshaft bearing surfaces on engine blocks.\nKnown casting processes can include quiescently feeding molten metal upwards into a casting chamber in a counter-gravity fill process. The casting process can include subsequently inverting the casting chamber to allow molten metal to gravity-feed into the inverted casting chamber to fully form the casting."}
{"text": "General Aspects and Inhibition of 17β-HSD1: Steroid hormones are important chemical carriers of information serving for the longterm and global control of cellular functions. They control the growth and the differentiation and function of many organs. On the other hand, they may also have negative effects and favor the pathogenesis and proliferation of diseases in the organism, such as mammary and prostate cancers (Deroo, B J. et al., J. Clin. Invest., 116: 561-570 (2006); Fernandez, S. V. et al., Int. J. Cancer, 118: 1862-1868 (2006)).\nThe biosynthesis of steroids takes place in the testes or ovaries, where sex hormones are produced. In addition, the production of glucocorticoids and mineral corticoids takes place in the adrenal glands. Moreover, individual synthetic steps also occur outside the glands, namely in the brain or in the peripheral tissue, e.g., adipose tissue (Bulun, S. E. et al., J. Steroid Biochem. Mol. Biol., 79: 19-25 (2001); Gangloff, A. et al., Biochem. J., 356: 269-276 (2001)). In this context, Labrie coined the term “intracrinology” in 1988 (Labrie, C. et al., Endocrinology, 123: 1412-1417 (1988); Labrie, F. et al., Ann. Endocrinol. (Paris), 56: 23-29 (1995); Labrie, F. et al., Horm. Res., 54: 218-229 (2000)). Attention was thus focused on the synthesis of steroids that are formed locally in peripheral tissues and also display their action there without getting into the blood circulation. The intensity of the activity of the hormones is modulated in the target tissue by means of various enzymes.\nThus, it could be shown that the 17β-hydroxysteroid dehydrogenase type 1 (17β-HSD1), which catalyzes the conversion of estrone (E1) to estradiol (E2), is more abundant in endometriotic tissue and breast cancer cells while there is a deficiency in 17β-HSD type 2, which catalyzes the reverse reaction (Bulun, S. E. et al., J. Steroid Biochem. Mol. Biol., 79: 19-25 (2001); Miyoshi, Y. et al., Int. J. Cancer, 94: 685-689 (2001)).\nA major class of steroid hormones is formed by the estrogens, the female sex hormones, whose biosynthesis takes place mainly in the ovaries and reaches its maximum immediately before ovulation. However, estrogens also occur in the adipose tissue, muscles and some tumors. Their main functions include a genital activity, i.e., the development and maintenance of the female sexual characteristics as well as an extragenital lipid-anabolic activity leading to the development of subcutaneous adipose tissue. In addition, they are involved in the pathogenesis and proliferation of estrogen-related diseases, such as endometriosis, endometrial carcinoma, adenomyosis, breast cancer and endometrial hyperplasia (Bulun, S. E. et al., J. Steroid Biochem. Mol. Biol., 79: 19-25 (2001); Miyoshi, Y. et al., Int. J. Cancer, 94: 685-689 (2001); Gunnarsson, C. et al., Cancer Res., 61: 8448-8451 (2001); Kitawaki, J., Journal of Steroid Biochemistry & Molecular Biology, 83: 149-155 (2003); Vihko, P. et al., J. Steroid. Biochem. Mol. Biol., 83: 119-122 (2002); Vihko, P. et al., Mol. Cell. Endocrinol., 215: 83-88 (2004); Saloniemi T. et al., Am. J. Pathol. 176: 1443-1451 (2010)), ovarian cancer, prostate cancer (Elo et al J. Cancer 66: 37 (1996)), acne (Odlind et al., Clin. Endocrinol. 16: 243-249 (1982)), androgen-dependent hair loss and psoriasis (Hughes et al. Endocrinology 138: 3711 (1997)).\nThe most potent estrogen is E2, which is formed in premenopausal females, mainly in the ovaries. On an endocrine route, it arrives at the target tissues, where it displays its action by means of an interaction with the estrogen receptor (ER) α. After the menopause, the plasma E2 level decreases to 1/10 of the E2 level found in premenopausal females (Santner, S. J. et al., J. Clin. Endocrinol. Metab., 59: 29-33 (1984)). E2 is mainly produced in the peripheral tissue, e.g., breast tissue, endometrium, adipose tissue and skin, from inactive precursors, such as estrone sulfate -EI-S), dehydroepiandrosterone (DHEA) and DHEA-S. These reactions occur with the participation of various steroidogenic enzymes (hydroxysteroid dehydrogenases, aromatase), which are in part more abundantly produced in the peripheral tissue, where these active estrogens display their action. As a consequence of such intracrine mechanism for the formation of E2-, its concentration in the peripheral tissue, especially in estrogen-related diseases, is higher than that in the healthy tissue. Above all, the growth of many breast cancer cell lines is stimulated by a locally increased E2 concentration. Further, the occurrence and progress of diseases such as endometriosis, leiomyosis, adenomyosis, menorrhagia, metrorrhagia and dysmenorrhea is dependent on a significantly increased -E2 level in accordingly diseased tissue.\nEndometriosis is an estrogen-related disease affecting about 5 to 10% of all females of childbearing age (Kitawaki, J., Journal of Steroid Biochemistry & Molecular Biology, 83: 149-155 (2003)). From 35 to 50% of the females suffering from abdominal pain and/or sterility show signs of endometriosis (Urdl, W., J. Reproduktionsmed. Endokrinol., 3: 24-30 (2006)). This disease is defined as a histologically detected ectopic endometrial glandular and stromal tissue. In correspondingly severe cases, this chronic disease, which tends to relapse, leads to pain of different intensities and variable character and possibly to sterility. Three macroscopic clinical pictures are distinguished: peritoneal endometriosis, retroperitoneal deep-infiltrating endometriosis including adenomyosis uteri, and cystic ovarial endometriosis. There are various explanatory theories for the pathogenesis of endometriosis, e.g., the metaplasia theory, the transplantation theory and the theory of autotraumatization of the uterus as established by Leyendecker (Leyendecker, G. et al., Hum. Reprod., 17: 2725-2736 (2002)).\nAccording to the metaplasia theory (Meyer, R., Zentralbl. Gynakol., 43: 745-750 (1919); Nap, A. W. et al., Best Pract. Res. Clin. Obstet. Gynaecol., 18: 233-244 (2004)), pluripotent coelomic epithelium is supposed to have the ability to differentiate and form endometriotic foci even in adults under certain conditions. This theory is supported by the observation that endometrioses, in part severe ones, can occur in females with lacking uterus and gynastresy. Even in males who were treated with high estrogen doses due to a prostate carcinoma, an endometriosis could be detected in singular cases.\nAccording to the theory postulated by Sampson (Halme, J. et al., Obstet. Gynecol., 64: 151-154 (1984); Sampson, J., Boston Med. Surg. J., 186: 445-473 (1922); Sampson, J., Am. J. Obstet. Gynecol., 14: 422-469 (1927)), retrograde menstruation results in the discharge of normal endometrial cells or fragments of the eutopic endometrium into the abdominal cavity with potential implantation of such cells in the peritoneal space and further development to form endometriotic foci. Retrograde menstruation could be detected as a physiological event. However, not all females with retrograde menstruation become ill with endometriosis, but various factors, such as cytokines, enzymes, growth factors (e.g., matrix metalloproteinases), play a critical role.\nThe enhanced autonomous non-cyclical estrogen production and activity as well as the reduced estrogen inactivation are typical peculiarities of endometriotic tissue. This enhanced local estrogen production and activity is caused by a significant overexpression of aromatase, expression of 17β-HSD1 and reduced inactivation of potent E2 due to a lack of 17β-HSD2, as compared to the normal endometrium (Bulun, S. E. et al., J. Steroid Biochem. Mol. Biol., 79: 19-25 (2001); Kitawaki, J., Journal of Steroid Biochemistry & Molecular Biology, 83: 149-155 (2003); Karaer, O. et al., Acta. Obstet. Gynecol. Scand., 83: 699-706 (2004); Zeitoun, K. et al., J. Clin. Endocrinol. Metab., 83: 4474-4480 (1998)).\nThe polymorphic symptoms caused by endometriosis include any pain symptoms in the minor pelvis, back pain, dyspareunia, dysuria and defecation complaints.\nOne of the therapeutic measures employed most frequently in endometriosis is the surgical removal of the endometriotic foci (Urdl, W., J. Reproduktionsmed. Endokrinol., 3: 24-30 (2006)). Despite new therapeutic concepts, medicamental treatment remains in need of improvement. The purely symptomatic treatment of dysmenorrhea is effected by means of non-steroidal anti-inflammatory drugs (NSAID), such as acetylsalicylic acid, indomethacine, ibuprofen and diclofenac. Since a COX2 overexpression could be observed both in malignant tumors and in the eutopic endometrium of females with endometriosis, a therapy with the selective COX2 inhibitors, such as celecoxib, suggests itself (Fagotti, A. et al., Hum. Reprod. 19: 393-397 (2004); Hayes, E. C. et al., Obstet. Gynecol. Surv., 57: 768-780 (2002)). Although they have a better gastro-intestinal side effect profile as compared to the NSAID, the risk of cardiovascular diseases, infarction and stroke is increases, especially for patients with a predamaged cardiovascular system (Dogne, J. M. et al., Curr. Pharm. Des., 12: 971-975 (2006)). The causal medicamental theory is based on estrogen deprivation with related variable side effects and a generally contraceptive character. The gestagens with their anti-estrogenic and antiproliferative effect on the endometrium have great therapeutic significance. The most frequently employed substances include medroxyprogesterone acetate, norethisterone, cyproterone acetate. The use of danazole is declining due to its androgenie side effect profile with potential gain of weight, hirsutism and acne. The treatment with GnRH analogues is of key importance in the treatment of endometriosis (Rice, V.; Ann. NY Acad. Sei., 955: 343-359 (2001)); however, the duration of the therapy should not exceed a period of 6 months since a longer term application is associated with irreversible damage and an increased risk of fracture. The side effect profile of the GnRH analogues includes hot flushes, amenorrhea, loss of libido and osteoporosis, the latter mainly within the scope of a long term treatment.\nAnother therapeutic approach involves the steroidal and non-steroidal aromatase inhibitors. It could be shown that the use of the non-steroidal aromatase inhibitor letrozole leads to a significant reduction of the frequency and severity of dysmenorrheal and dyspareunia and to a reduction of the endometriosis marker CA125 level (Soysal, S. et al., Hum. Reprod., 19: 160-167 (2004)). The side effect profile of aromatase inhibitors ranges from hot flushes, nausea, fatigue to osteoporosis and cardiac diseases. Long term effects cannot be excluded. All the possible therapies mentioned herein are also employed in the combatting of diseases such as leiomyosis, adenomyosis, menorrhagia, metrorrhagia and dysmenorrhea.\nEvery fourth cancer disease in the female population falls under the category of mammary cancers. This disease is the main cause of death in the western female population at the age of from 35 to 54 years (Nicholls, P. J., Pharm. J., 259: 459-470 (1997)). Many of these tumors exhibit an estrogen-dependent growth and are referred to as so-called HDSC (hormone dependent breast cancer). A distinction is made between ER+ and ER− tumors. The classification criteria are important to the choice of a suitable therapy. About 50% of the breast cancer cases in premenopausal females and 75% of the breast cancer cases in postmenopausal females are ER+ (Coulson, C., Steroid biosynthesis and action, 2nd edition, 95-122 (1994); Lower, E. et al., Breast Cancer Res. Treat., 58: 205-211 (1999)), i.e., the growth of the tumor is promoted by as low as physiological concentrations of estrogens in the diseased tissue.\nThe therapy of choice at an early stage of breast cancer is surgical measures, if possible, breast-preserving surgery. Only in a minor number of cases, mastectomy is performed. In order to avoid relapses, the surgery is followed by radiotherapy, or else radiotherapy is performed first in order to reduce a larger tumor to an operable size. In an advanced state, or when metastases occur in the lymph nodes, skin or brain, the objective is no longer to heal the disease, but to achieve a palliative control thereof.\nThe therapy of the mammary carcinoma is dependent on the hormone receptor status of the tumor, on the patient's hormone status and on the status of the tumor (Paepke, S. et al., Onkologie, 26 Suppl., 7: 4-10 (2003)). Various therapeutical approaches are available, but all are based on hormone deprivation (deprivation of growth-promoting endogenous hormones) or hormone interference (supply of exogenous hormones). However, a precondition of such responsiveness is the endocrine sensitivity of the tumors, which exists with HDSC ER+ tumors. The drugs employed in endocrine therapy include GnRH analogues, anti-estrogens and aromatase inhibitors. GnRH analogues, such as gosereline, will bind to specific membrane receptors in the target organ, the pituitary gland, which results in an increased secretion of FSH and LH. These two hormones in turn lead to a reduction of GnRH receptors in a negative feedback loop in the pituitary cells. The resulting desensitization of the pituitary cells towards GnRH leads to an inhibition of FSH and LH secretion, so that the steroid hormone feedback loop is interrupted. The side effects of such therapeutic agents include hot flushes, sweats and osteoporosis.\nAnother therapeutic option is the use of anti-estrogens, antagonists at the estrogen receptor. Their activity is based on the ability to competitively bind to the ER and thus avoid the specific binding of the endogenous estrogen. Thus, the natural hormone is no longer able to promote tumor growth. Today, therapeutic use involves so-called SERM (selective estrogen receptor modulators), which develop estrogen agonism in tissues such as bones or liver, but have antagonistic andjor minimal agonistic effects in breast tissue or uterus (Holzgrabe, U., Pharm. Unserer Zeit, 33: 357-359 (2004); Pasqualini, J. R., Biochim. Biophys. Acta., 1654: 123-143 (2004); Sexton, M. J. et al., Prim Care Update Ob Gyns, 8: 25-30 (2001)). Thus, these compounds are not only effective in combatting breast cancer, but also increase the bone density and reduce the risk of osteoporosis in postmenopausal females. The use of the SERM tamoxifen is most widely spread. However, after about 12-18 months of treatment, there is development of resistance, an increased risk of endometrial cancers and thrombo-embolic diseases due to the partial agonistic activity at the ER (Goss, P. E. et al., Clin. Cancer Res., 10: 5717-5723 (2004); Nunez, N. P. et al., Clin. Cancer Res., 10: 5375-5380 (2004)).\nThe enzymatically catalyzed estrogen biosynthesis may also be influenced by selective enzyme inhibitors. The enzyme aromatase, which converts C19 steroids to C18 steroids, was one of the first targets for lowering the E2 level. This enzyme complex, which belongs to the cytochrome P-450 enzymes, catalyzes the aromaticzation of the androgenic A ring to form estrogens. The methyl group at position 10 of the steroid is thereby cleaved off. The first aromatase inhibitor employed for the therapy of breast cancer was aminoglutethimide. However, aminoglutethimide affects several enzymes of the cytochrome P-450 superfamily and thus inhibits a number of other biochemical conversions. For example, among others, the compound interferes with the steroid production of the adrenal glands so heavily that a substitution of both glucocorticoids and mineralocorticoids may be necessary. In the meantime, more potent and more selective aromatase inhibitors, which can be subdivided into steroidal and non-steroidal compounds, are on the market. The steroidal inhibitors include, for example, exemestane, which has a positive effect on the bone density, which is associated with its affinity for the androgen receptor (Goss, P. E. et al., Clin. Cancer Res., 10: 5717-5723 (2004)). However, this type of compounds is an irreversible inhibitor that also has a substantial number of side effects, such as hot flushes, nausea, fatigue. However, there are also non steroidal compounds that are employed therapeutically, for example, letrozole. The advantage of these compounds resides in the lesser side effects, they do not cause uterine hypertrophy, but have no positive effect on the bone density and result in an increase of LDL (Iow density lipoprotein), cholesterol and triglyceride levels (Goss, P. E. et al., Clin. Cancer Res., 10: 5717-5723 (2004); Nunez, N. P. et al., Clin. Cancer Res., 10: 5375-5380 (2004)). Today, aromatase inhibitors are predominantly employed as second-line therapeutic agents. In the meantime, however, the equivalence or even superiority of aromatase inhibitors to SERM, such as tamoxifene, has been proven in clinical studies (Geisler, J. et al., Crit. Rev. Oncol. Hematol., 57: 53-61 (2006); Howell, A. et al., Lancet, 365: 60-62 (2005)). Thus, the use of aromatase inhibitors also as first-line therapeutical agents is substantiated.\nHowever, the estrogen biosynthesis in the peripheral tissue also includes other pathways for the production of E1 and the more potent E2 by avoiding the enzyme aromatase that is locally present in the target tissue, for example, breast tumors. Two pathways for the production of estrogens in breast cancer tissue are postulated (Pasqualini, J. R., Biochim. Biophys. Acta., 1654: 123-143 (2004)), the aromatase pathway (Abul-Hajj, Y J. et al., Steroids, 33: 205-222 (1979); Lipton, A. et al., Cancer, 59: 779-782 (1987)) and the sulfatase pathway (Perei, E. et al., J. Steroid. Biochem., 29: 393-399 (1988)). The aromatase pathway includes the production of estrogens from androgens with participation of the enzyme aromatase. The sulfatase pathway is the pathway for the production of E1/E2 by means of the enzyme steroid sulfatase, an enzyme that catalyzes the conversion of E1 sulfate and DHEA-S to estrone and DHEA. In this way, 10 times as much E1 is formed in the target tissue as compared to the aromatase pathway (Santner, S. J. et al., J. Clin. Endocrinol. Metab., 59: 29-33 (1984)). E1 is then reduced by means of the enzyme 17β-HSD1 to form E2, the most potent estrogen. Steroid sulfatase and 17β-HSD1 are new targets in the battle against estrogen-related diseases, especially for the development of therapeutic agents for mammary carcinomas (Pasqualini, J. R., Biochim. Biophys. Acta., 1654: 123-143 (2004)).\nNumerous steroidal sulfatase inhibitors could be found, including the potent irreversible inhibitor EMATE, which exhibited an agonistic activity at the estrogen receptor, however (Ciobanu, L. C. et al., Cancer Res., 63: 6442-6446 (2003); Hanson, S. R. et al., Angew. Chem. Int. Ed. Engl., 43: 5736-5763 (2004)). Some potent non-steroidal sulfatase inhibitors could also be found, such as COUMATE and derivatives as weil as numerous sulfamate derivatives of tetrahydronaphthalene, indanone and tetralone (Hanson, S. R. et al., Angew. Chem. Int. Ed. Engl., 43: 5736-5763 (2004)). Recently, one sulfatase inhibitor has completed a phase I clinical trial in postmenopausal women with breast cancer (Foster, P. A. et al., Anticancer Agents Med Chem. 8(7):732-8 (2008).\nThe inhibition of 17β-HSD1, a key enzyme in the biosynthesis of E2, the most potent estrogen, could suggest itself as an option in the therapy of estrogen-related diseases in both premenopausal and postmenopausal females (Kitawaki, J., Journal of Steroid Biochemistry & Molecular Biology, 83: 149-155 (2003); Allan, G. M. et al., Mol. Cell Endocrinol., 248: 204-207 (2006); Penning, T. M., Endocr. Rev., 18: 281-305 (1997); Sawicki, M. W. et al., Proc. Natl. Acad. Sci. USA, 96: 840-845 (1999); Vihko, P. et al., Mol. Cell. Endocrinol., 171: 71-76 (2001)). An advantage of this approach is the fact that the intervention is effected in the last step of estrogen biosynthesis, i.e., the conversion of E1 to the highly potent E2 is inhibited. The intervention is effected in the biosynthetic step occurring in the peripheral tissue, so that a reduction of E2 production takes place locally in the diseased tissue. The use of correspondingly selective inhibitors would probably be associated with little side effects since the synthesis of other steroids would remain unaffected. To achieve a selective effect, it would be important that such inhibitors exhibit no or only very little agonistic activity at the ER, especially at the ER α, since agonistic binding is accompanied by an activation and thus proliferation and differentiation of the target cell. In contrast, an antagonistic activity of such compounds at the ER would be tolerated since the inhibitor would prevent the natural substrates from binding at the receptor which in turn will result in a further reduction of the proliferation of the target cells. The use of selective 17β-HSD1 inhibitors for the therapy of numerous estrogen-dependent diseases is discussed, for example, for breast cancer, tumors of the ovaries, prostate carcinoma, endometrial carcinoma, endometriosis, adenomyosis, endometrial hyperplasia, acne, psoriasis, and androgen-dependent hair loss. Highly interesting and completely novel is the proposal to employ selective inhibitors of 17β-HSD1 for prevention when there is a genetic disposition for breast cancer (Miettinen, M. et al., J. Mammary Gland. Biol. Neoplasia, 5: 259-270 (2000)).\nHydroxysteroid dehydrogenases (HSD) can be subdivided into different classes. The 11β-HSD modulate the activity of glucocorticoids, 3β-HSD catalyzes the reaction of Δ5-3β-hydroxysteroids (DHEA or 5-androstene-3β,17β-diol) to form Δ5-3β-ketosteroids (androstenedione or testosterone). 17β-HSDs convert the less active 17-ketosteroids to the corresponding highly active 17-hydroxy compounds (androstenedione to testosterone and E1 to E2) or conversely (Payne, A. H. et al., Endocr. Rev., 25: 947-970 (2004); Peltoketo, H. et al., J. Mol. Endocrinol., 23: 1-11 (1999); Suzuki, T. et al., Endocr. Relat. Cancer, 12: 701-720 (2005)). Thus, the HSD play a critical role in both the activation and the inactivation of steroid hormones. Depending on the cell's need for steroid hormones, they alter the potency of the sex hormones (Penning, T. M., Endocr. Rev., 18: 281-305 (1997)), for example, Ei is converted to the highly potent E2 by means of 17β-HSD1, while E2 is converted to the less potent E1 by means of 17β-HSD2; 17β-HSD2 inactivates E2 while 17β-HSD1 activates E1.\nTo date, fourteen different mammalian 17β-HSDs have been identified (Haller, F. et al., J. Mol. Biol. doi: 10.1016jj.jmb.2010.04.002 (2010); Zhongyi, S. et al., Endocrinology 148 3827-3836 (2007); Miyoshi, T. et al., Int. J. Cancer 94 (2001) 685-689), and twelve of these enzymes could be cloned (Suzuki, T. et al., Endocr. Relat. Cancer, 12: 701-720 (2005)). They all belong to the so-called short chain dehydrogenasejreductase (SDR) family, with the exception of 17β-HSD5, which is a ketoreductase. The amino acid identity between the different 17β-HSDs is as low as 20-30% (Luu-The, V., J. Steroid Biochem. Mol. Biol., 76: 143-151 (2001)), and they are membrane-bound or soluble enzymes. The X-ray structure of 6 human subtypes is known (Ghosh, D. et al., Structure, 3: 503-513 (1995); Kissinger, C. R. et al., J. Mol. Biol., 342: 943-952 (2004); Zhou, M. et al., Acta Crystallogr. D. Biol. Crystallogr., 58: 1048-1050 (2002). The 17β-HSDs are NAD(H)-dependent and NADP(H)-dependent enzymes. They play a critical role in the hormonal regulation in humans. The enzymes are distinguished by their tissue distribution, catalytic preference (oxidation or reduction), substrate specificity and subcellular localization. The same HSD subtype was found in different tissues. It is likely that all 17β-HSDs are expressed in the different estrogen-dependent tissues, but in different concentrations. In diseased tissue, the ratio between the different subtypes is altered as compared to healthy tissue, some subtypes being overexpressed while others may be absent. This may cause an increase or decrease of the concentration of the corresponding steroid. Thus, the 17β-HSDs play an extremely important role in the regulation of the activity of the sex hormones. Further, they are involved in the development of estrogen-sensitive diseases, such as breast cancer, ovarian, uterine and endometrial carcinomas, as weil as androgenrelated diseases, such as prostate carcinoma, benign prostate hyperplasia, acne, hirsutism etc. It has been shown that some 17β-HSDs are also involved in the development of further diseases, e.g., pseudohermaphrodism (17β-HSD3 (Geissler, W. M. et al., Nat. Genet., 7: 34-39 (1994))), bifunctional enzyme deficiency (17β-HSD4 (van Grunsven, E. G. et al., Proc. Natl. Acad. Sci. USA, 95: 2128-2133 (1998))), polycystic kidney diseases (17β-HSD8 (Maxwell, M. M. et al., J. Biol. Chem., 270: 25213-25219 (1995))) and Alzheimer's (17β-HSD10 (Kissinger, C. R. et al., J. Mol. Biol., 342: 943-952 (2004); He, X. Y. et al., J. Biol. Chem., 274: 15014-15019 (1999); He, X. Y. et al., Mol. Cell Endocrinol., 229: 111-117 (2005); He, X. Y. et al., J. Steroid Biochem. Mol. Biol., 87: 191-198 (2003); Yan, S. D. et al., Nature, 389: 689-695 (1997))). The best characterized member of the 17β-HSDs is the type 1 17β-HSD. The 17β-HSD1 is an enzyme from the SDR family, also referred to as human placenta E2 dehydrogenase (Gangloff, A. et al., Biochem. J., 356 269-276 (2001); Jomvall, H. et al., Biochemistry, 34 6003-6013 (1995)). Its designation as assigned by the enzyme commission is E.C.1.1.1.62.\nEngel et al. (Langer, L. J. et al., J. Biol. Chem., 233: 583-588 (1958)) were the first to describe this enzyme in the 1950's. In the 1990's, first crystallization attempts were made, so that a total of 20 crystallographic structures can be recurred to today in the development of inhibitors (Negri, M. et al. PLoS ONE 5(8): e12026. doi: 10.1371/journal.pone.0012026 (2010)). Available are X-ray structures of the enzyme alone, but also of binary and ternary complexes of the enzyme with its substrate and other ligands or substrate/ligand and cofactor.\n17β-HSD1 is a soluble cytosolic enzyme. NADPH serves as a cofactor. 17β-HSD1 is encoded by a 3.2 kb gene consisting of 6 exons and 5 introns that is converted to a 2.2 kb transcript (Luu-The, V., J. Steroid Biochem. Mol. Biol., 76: 143-151 (2001); Labrie, F. et al., J. Mol. Endocrinol., 25: 1-16 (2000)). It is constituted by 327 amino acids. The molecular weight of the monomer is 34.9 kDa (Penning, T. M., Endocr. Rev., 18: 281-305 (1997)). 17β-HSD1 is expressed in the placenta, liver, ovaries, endometrium, prostate gland, peripheral tissue, such as adipose tissue and breast cancer cells (Penning, T. M., Endocr. Rev., 18: 281-305 (1997)). It was isolated for the first time from human placenta (Jarabak, J. et al., J. Biol. Chem., 237: 345-357 (1962)). The main function of 17β-HSD1 is the conversion of the less active E1 to the highly potent E2. However, it also catalyzes to a lesser extent the reaction of dehydroepiandrosterone (DHEA) to 5-androstene-3β,17β-diol, an androgen showing estrogenic activity (Labrie, F., Mol. Cell. Endocrinol., 78: CI13-118 (1991); Poirier, D., Curr. Med. Chem., 10: 453-477 (2003); Poulin, R. et al., Cancer Res., 46: 4933-4937 (1986)). In vitro, the enzyme catalyzes the reduction and oxidation between E1 and E2 while it catalyzes only the reduction under physiological conditions. These bisubstrate reactions proceed according to a random catalytic mechanism, i.e., either the steroid or the cofactor is first to bind to the enzyme (Negri, et al. PLoS ONE 5(8): e12026. doi:10.1371/journal.pone.0012026 (2010)). The enzyme consists of a substrate binding site and a channel that opens into the cofactor binding site. The substrate binding site is a hydrophobie tunnel having a high complementarity to the steroid. The 3-hydroxy and 17-hydroxy groups in the steroid form four hydrogen bonds to the amino acid residues His221, Glu282, Ser142 and Tyr155. The hydrophobie van der Waals interactions seem to form the main interactions with the steroid while the hydrogen bonds are responsible for the specificity of the steroid for the enzyme (Labrie, F. et al., Steroids, 62: 148-158 (1997)). Like with all the other enzymes of this family, what is present as a cofactor binding site is the Rossmann fold, which is a region consisting of ex-helices and β-sheets (β-α-β-α-β)2, a generally occurring motif Gly-Xaa-Xaa-Xaa-Gly-Xaa-Gly, and a nonsense region Tyr-Xaa-Xaa-Xaa-Lys within the active site. What is important to the activity is a catalytic tetrade consisting of Tyr155-Lys159-Ser142-Asnl14, which stabilize the steroid and the ribose in the nicotinamide during the hydride transfer (Alho-Richmond, S. et al., Mol. Cell Endocrinol., 248: 208-213 (2006); Labrie, F. et al., Steroids, 62: 148-158 (1997); Nahoum, V. et al., Faseb. J., 17: 1334-1336 (2003)).\nThe gene encoding 17β-HSD1 is linked with the gene for mammary and ovarian carcinomas that is very susceptible to mutations and can be inherited, the BRCA1 gene, on chromosome 17q11-q21 (Labrie, F. et al., J. Mol. Endocrinol., 25: 1-16 (2000)). As has been demonstrated, the activity of 17β-HSD1 is higher in endometrial tissue and breast cancer cells as compared to healthy tissue, which entails high intracellular E2 levels, which in turn cause proliferation and differentiation of the diseased tissue (Bulun, S. E. et al., J. Steroid Biochem. Mol. Biol., 79: 19-25 (2001); Miyoshi, Y. et al., Int. J. Cancer, 94: 685-689 (2001); Kitawaki, J., Journal of Steroid Biochemistry & Molecular Biology, 83: 149-155 (2003); Pasqualini, J. R., Biochim. Biophys. Acta., 1654: 123-143 (2004); Vihko, P. et al., Mol. Cell. Endocrinol., 171: 71-76 (2001); Miettinen, M. et al., Breast Cancer Res. Treat., 57: 175-182 (1999); Sasano, H. et al., J. Clin. Endocrinol. Metab., 81: 4042-4046 (1996); Yoshimura, N. et al., Breast Cancer Res., 6: R46-55 (2004)). An inhibition of 17β-HSD1 could result in the E2 level being lowered and thus lead to a regression of the estrogen-related diseases. Further, selective inhibitors of 17β-HSD1 could be used for prevention when there is a genetic disposition for breast cancer (Miettinen, M. et al., J. Mammary Gland. Biol. Neoplasia, 5: 259-270 (2000)).\nTherefore, this enzyme would suggest itself as a target for the development of novel selective and non-steroidal inhibitors as therapeutic agents in the battle against estrogen-related diseases. Recently, in vivo efficacy of 17β-HSD1 inhibitors has been reported in two animal models. Immunodeficient mice were inoculated either with MCF-7 cells over-expressing human recombinant 17β-HSD1 enzyme (Husen, B. et al., Mol. Cell. Endocrinol., 248, 109-113 (2006); Husen, B. et al., Endocrinology, 147, 5333-5339 (2006)) or with T47D cells naturally expressing 17β-HSD1 (Day, J. M. et al.; Int. J. Cancer, 122, 1931-1940 (2008). In both models, the El induced tumor growth was reduced by 17β-HSD1 inhibitors, validating 17β-HSD1 as a novel target for the treatment of estrogen dependent diseases. Up to date however, no 17β-HSD1 inhibitor has entered clinical trials.\nIn the literature, only a few compounds have been described as inhibitors of 17βHSD1 (D. Poirier, Anticancer Agents Med. Chem. 9 642-660 (2009); Day, J. M. et al., Minerva Endocrinol., 35(2), 87-108 (2010); D. Poirier, Expert Opin Ther Pat., doi:10.1517/13543776.2010.505604 (2010)). Most inhibitors are steroidal compounds obtained by different variations of the estrogen skeleton (Rouillard, F. et al., Open Enzyme Inhib. J., 1 61-71 (2008); D. Poirier, Anticancer Agents Med. Chem. 9 642-660 (2009); Mazumdar M. et al.; Biochem J., 424(3):357-366 (2009); Möller G et al. Bioorg Med Chem Lett., 19(23): 6740-6744 (2009); Berube M. et al. Molecules 15 1590-1631 (2010)).\n\nAnother class of compounds which has been described is the so-called hybrid inhibitors (Berube, M. et al., Can. J. Chem. 87 1180-1199 (2009)), compounds that, due to their molecular structure, not only attack at the substrate binding site, but also undergo interactions with the cofactor binding site. The inhibitors have the following structure:                adenosine moiety or simplified derivatives that can interact with the cofactor binding site;        E2 or E1 moiety, which interacts with the substrate binding site; and        a spacer of varying length as a linking element between the two moieties.        \n\nAmong these compounds, inhibitors have been synthesized that exhibit a good inhibition of the enzyme and a good selectivity for 17β-HSD2 (compound B (Lawrence, H. R. et al., J. Med. Chem., 48: 2759-2762 (2005)). In addition, the inventors consider that a small estrogenic effect can be achieved by a substitution at the C2 of the steroid skeleton (Cushman, M. et al., J. Med. Chem., 38: 2041-2049 (1995); Leese, M. P. et al., J. Med. Chem., 48: 5243-5256 (2005)); however, this effect has not yet been demonstrated in tests.\nHowever, a drawback of these steroidal compounds may be a low selectivity. With steroids, there is a risk that the compounds will also interfere with other enzymes of the steroid biosynthesis, which would lead to side effects. In addition, due to their steroidal structure, they may have an affinity for steroid receptors and function as agonists or antagonists.\nAmong the phytoestrogens, which have affinity for the estrogen receptor and act as estrogens or anti-estrogens depending on the physiological conditions, flavones, isoflavones and lignans have been tested for an inhibitory activity (Makela, S. et al., Proc. Soc. Exp. Biol. Med., 217: 310-316 (1998); Makela, S. et aL, Proc. Soc. Exp. Biol. Med., 208: 51-59 (1995); Brooks, J. D. et al., J. Steroid Biochem. Mol. Biol., 94: 461-467 (2005)). Coumestrol was found to be particularly potent, but of course showed estrogenic activity (Nogowski, L., J. Nutr. Biochem., 10: 664-669 (1999)). Gossypol derivatives were also synthesized as inhibitors (US2005/0228038). In this case, however, the cofactor binding site rather than the substrate binding site is chosen as the target site (Brown, W. M. et al., Chem. Biol. Interact., 143-144, 481-491 (2003)), which might entail problems in selectivity with respect to other enzymes utilizing NAD(H) or NADP(H).\n\nIn addition to diketones, such as 2,3-butanedione and glyoxal, which were used for studies on the enzyme, suicide inhibitors were also tested. However, these were found not to be therapeutically utilizable since the oxidation rate of the alcohols to the corresponding reactive form, namely the ketones, was too weak (Poirier, D., Curr. Med. Chem., 10: 453-477 (2003)).\nIn other studies, Jarabak et al. (Jarabak, J. et al., Biochemistry, 8: 2203-2212 (1969)) examined various non-steroidal inhibitors for their inhibitory effect, U-11-100A having been found as the most potent compound in this group. However, as compared to other non steroidal compounds, U-II-100A is a weak inhibitor of 17βHSD1.\n\nRecently, using a pharmacophore model for 17β-HSD1 inhibitors Schuster et al. identified some compounds with 17β-HSD1 inhibitory activity in the micromolar range (Schuster, D. et al., J Med Chem. 51:4188-4199 (2008)). Regarding additional non steroidal 17β-HSD1 inhibitors, 5 templates reveal interesting biological activities: A. Thiophenepyrimidinones (US2005/038053; Messinger, J. et al., Mol. Cell. Endocrinol., 248: 192-198 (2006); W02004jll0459; Lilienkampf, A. et al., J. Med. Chem. 52: 6660-6671 (2009)); B. Biphenyl ethanones (Allan, G. M. et al. Bioorg. Med. Chem. 16: 4438-4456 (2008)); C. Hydroxyphenylnaphtols (WO/08EP/53672; Marchais-Oberwinkler S. et al., J. Med. Chem., 51: 4685-4698 (2008)) Marchais-Oberwinkler, S. et al., J. Med. Chem., 54: 534-547, 2011); D. Heterocyclic substituted biphenylols (Oster, A. et al., Bioorg. Med. Chem., 18: 3494-3505 (2010)); E. Bis(hydroxyphenyl) arenes (W02009/02746; Bey, E. et al., J. Med. Chem., 52: 6724-6743 (2009)).\n\nMost of those classes showed high potency at the protein level (IC50<20 nM; W02004/53424) but a limited inhibitory activity (IC50>200 nM) in cell-based 17β-HSD1 assays (Messinger J. et al, Mol. Cell. Endocrinol., 248: 192-198 (2006) and Bey E. et al, J. Med. Chem., 52: 6724-6743 (2009)), which might be due to poor cell membrane permeability.\nRecently described coumarins display only moderate inhibition of 17β-HSD1 (Star{hacek over (c)}ević et al, J. Med. Chem. 54: 248-261 (2011)), whereas bicyclic substituted hydroxyphenylmethanones have been described as potent inhibitors (Oster et al., J. Med. Chem. 53: 8176-8186 (2010)).\n\nInhibitors with hydroxybenzothiazole core (Spadaro et al. PLoS ONE 7: e29252. doi:10.1371/journal.pone.0029252; J. Med. Chem. 55: 2469-2473 (2012)) are potent but show poor pharmacokinetics. Further, 2-benzoylbenzothiazole derivatives having lipid lowering activity are known from EP-A-735029, and N-(benzothiazol-2-yl)arylcarboxamide and I-(benzothiazol-2-yl)-3-(aryl)urea derivatives and their use for the inhibition of ubiquitination are known from WO2005/037845. Similar N-(benzothiazol-2-yl)arylcarboxamide and 1-(benzothiazol-2-yl)-3-(aryl)urea derivatives including frentizole (1-(6-methoxy-1,3-benzothiazol-2-yl)-3-phenylurea) are also known to interact with amyloid beta (Aβ) peptide and/or (Aβ)-binding alcohol dehydrogenase and are potential anti-Alzheimer agent (Xie, Y. et al., Bioorg. & Med. Chem. Lett. 16: 4657-60 (2006)).\nFinally, selective 17β-HSD1 inhibitors with (phenylthiazolyl)(phenyl)methanone (phenylthienyl)(phenyl)methanone and (benzothiazolyl)(phenyl)methanone structure are known from WO2012/025638, and selective 17β-HSD1 inhibitors with (phenyl-1,3-thiazol-4-yl)phenol and (phenylthienyl)phenol structure are known from DE102007040243A1.\n17Beta-Hydroxysteroid Dehydrogenase Type 2 Inhibitors: Estrogens and androgens play a crucial role in the development, growth and function of all tissues involved in reproduction and fertility. It is also well known, that E2 and testosterone/dihydrotestosterone (T/DHT) the most active estrogen and androgen, respectively, can be involved in a series of hormone-sensitive diseases. For example estrogens or androgens are often responsible for the development of breast cancer or prostate cancer, respectively, via stimulation of cell proliferation in the corresponding tissues (Travis, R. C. et al., Breast Cancer Res., 5: 239-247 (2003); Wilding, G., Cancer Surv., 14113-14130 (1992)) or insufficient levels of E2 and T/DHT predispose the human skeleton to osteoporosis in both men and women (Pietschmann, P. et al., Gerontology, 55:3-12 (2008)).\nOsteoporosis is a systemic skeletal disease characterized by deterioration of bone tissue and low bone mass, resulting in increased fragility of the bone and higher risk of fractures of the hip, spine and wrist. Osteoporotic fractures lead to pain, occasional disability and more important, they increase mortality (Cree, M. et al., J. Am. Geriatr. Soc., 48:283-288 (2000)).\nIn healthy individuals, bone mass is maintained by a balance between bone resorption and bone formation performed by the osteoclasts (OCs) and osteoblasts (OBs), respectively. This process of bone remodelling facilitates repair of microdamage, provides calcium uptake and therefore brings stability and strength to the bone.\nBone loss is, however, accelerated in post-menopausal women and in elderly men. The mechanisms by which elderly people, both men and women, lose bone are not fully understood and remain under investigation. Decreased quantity of sex hormones is one important factor causing bone loss. In women at menopause, estrogen deficiency (Cree, M. et al., J. Am. Geriatr. Soc., 48:283-288 (2000)) and in older men, estrogen and androgen insufficiency (Fink, H. A. et al., J. Clin. Endocrinol. Metab., 91:3908-3915 (2006); Meier, C. et al., Arch. Intern. Med., 168:47-54 (2008)) result in a disproportionate increase in bone loss as compared with bone formation and often lead to osteoporosis.\nSince OBs and OCs express estrogen receptors (Hoyland, J. A. et al., Bone, 20:87-92 (1997)) and respond to estrogen treatments, the most potent estrogen, E2 must have a direct effect on maintenance of bone mineral density. There are also substantial evidence that androgens like testosterone (T) and dihydrotestosterone (DHT) may as well be involved in bone formation, increasing OB activity and therefore protecting the bones against osteoporosis (Vanderschueren, D. et al., Endocr. Rev., 25:389-425 (2004), Vanderschueren, D. et al., Curr. Opin. Endocrinol. Diabetes Obes., 15:250-254 (2008)). T and DHT might act via activation of the androgen receptor (AR) which is present in the OBs (Bland, R., Clin. Sci., 98:217-240 (2000)) or could be the precursor of estrogens (the enzyme aromatase, responsible for the transformation of androgens in estrogens has been identified in the OBs). Controlled increase of E2 and T in bones of osteoporotic patients will simultaneously lower bone resorption (effect of E2) and raise bone formation (effect of T) improving bone loss and osteoporotic fractures. Augmentation of E2 and T levels in bones might be achieved by inhibition of the enzyme 17β-hydroxysteroid dehydrogenase type 2 (17β-HSD2) which has been identified in OBs (Dong, Y. et al., J. Bone Miner. Res., 13:1539-1546 (1998); Feix, M. et al., Mol. Cell. Endocrinol., 171:163-164 (2001); Janssen, J. Met al., J. Cell. Biochem., 75:528-37 (1999)). A restricted increase in E2 and T in OCs and OBs is important to avoid unwanted side-effects such as induction of breast cancer or prostate cancer. This might be achieved via an intracrine mechanism (Labrie, F., Mol. Cell. Endocrinol. 78:C113-118 (1991)), i.e., the transformation of E2 and T in inactive precursors should be blocked dominantly in the bone cells 17β-HSD2 (Wu, L. et al., J. Biol. Chem., 268:12964-12969 (1993); Lu, M., J. Biol. Chem., 277:22123-22130 (2002)) belongs to the hydroxysteroid dehydrogenase (HSD) superfamily. The HSD enzymes play pivotal roles not only in the activation but also in the inactivation of steroid hormones. Depending on the need of the cell, they modulate the potency of the sex hormones (Penning, T. M. et al., Endocr. Rev., 18:281-305 (1997)). For example, 17β-HSD1 activates E1 into the very potent E2, while 17β-HSD2 oxidizes E2 into E1 thus decreasing the action of the potent E2. It is therefore believed that in physiology 17β-HSD2 protects the cell against excessive level of active estrogen. To date, fourteen different 17β-HSDs have been identified (Lukacik, P. et al., Mol. Cell. Endocrinol. 248:61-71 (2006); Luu-The, V. et al., Best Pract. Res. Clin. Endocrinol. Metab., 22:207-221 (2008)) and twelve different subtypes have been cloned from human tissues. They all belong to the Short-Chain Dehydrogenase/Reductase family (except type 5). The 17β-HSDs show little amino acid identity (20-30%) and are membrane-bound or soluble enzymes. The X-ray structures of six human subtypes (type 1, 4, 8, 10, 11, 14) are known (Moeller, G. et al., Mol. Cell. Endocrinol. 301:7-19 (2009)). All of them are NAD(H)- or NADP(H)-dependent. The 17β-HSDs play therefore a key role in hormonal regulation and function in humans. They differ in tissue distribution, catalytic preference (oxidation or reduction), substrate specificity and subcellular localisation. It is likely that all 17β-HSDs, modulating estrogen and androgen action, are expressed in the different estrogen/androgen-dependent tissues but certainly at different concentrations. In diseased tissues, the ratio between the different subtypes is changed compared to healthy tissues: one enzyme subtype might be overexpressed or completely absent, leading to higher/lower concentrations of the specific steroids. Selective inhibition of one subtype could therefore be a good strategy to influence the level in estrogen and androgen in hormone sensitive diseases. 17β-HSD2 (EC 1.1.1.51) is a transmembrane protein localized in the endoplasmic reticulum (Wu, L. et al., J. Biol. Chem., 268:12964-12969 (1993); Lu, M., J. Biol. Chem., 277:22123-22130 (2002)). Its 3D-structure remains unknown. 17β-HSD2 shows a broad tissue distribution in adult: it is expressed in liver, small intestine, endometrium, urinary tracts and bone osteoblastic (Dong, Y. et al., J. Bone Miner. Res., 13:1539-1546 (1998); Feix, M. et al., Mol. Cell. Endocrinol., 171:163-164 (2001); Eyre, L. J. et al., J. Bone Miner. Res., 13:996-1004 (1998)) and osteoclastic (van der Eerden, B. C. J. et al., J. Endocrinol., 180:457-467 (2004)) cells. It has a predominant oxidative activity on hydroxy groups at the position C17 of the steroids: it converts the highly active 17β-hydroxysteroid estrogen E2 as well as the 17β-hydroxysteroid androgens T and DHT into their inactive keto forms using the cofactor NAD+. The broad tissue distribution together with the oxidative activity of 17β-HSD2 suggests that this enzyme may play an essential role in the inactivation of highly active estrogens and androgens, resulting in diminished sex hormone action in target tissues.\nRelatively few inhibitors of 17β-HSD2 have been described. Sam, K. M., J. Med. Chem., 38:4518-4528 (1995); Poirier, D. et al., Mol. Cell. Endocrinol., 171:119-128 (2001); Bydal, P. et al., Eur. J. Med. Chem., 44:632-644 (2009)) reported on the potent steroidal C17-spirolactone 1 (inhibitory activity of 70% at 1 μM and 40% at 100 nM). Cook et al. (Cook, J. H. et al., Tetrahedron Letters, 46:1525-1528 (2005); Gunn, D., Bioorg. Med. Chem. Lett., 15:3053-3057 (2005); Wood, J. et al., Bioorg. Med. Chem. Lett., 16:4965-4968 (2006)) found, using high through-put screening methods, the first new class of potent non-steroidal inhibitors of 17β-HSD2: the 4,5-disubstituted cis-pyrrolidines. The most active compound described is 2 (IC50=10 nM). It has been recently proven in monkeys, using an osteoporosis model for in vivo evaluation of the 17β-HSD2 inhibitor 3, that inhibition of this enzyme helps to maintain a sufficient local level of E2 and T in bones when the level of circulating active sex steroid drops (Bagi, C. M. et al., J. Musculoskelet. Neuronal Interact., 8:267-280 (2008)).\n\nThis modulation of steroid levels might be useful for a variety of indications like prevention and treatment of diseases caused by a misbalance of OB/OC activity like osteoporosis, osteopenia and impaired bone fracture healing (Undsay, R. et al., Obst. Gynecol., 76:290-295 (1990); Turner, R. T. et al., Endocr. Rev., 16:275-300 (1994); Meczekalski, B. et al., Gynecol. Endocrinol., 26:652-657 (2010)), postmenopausal symptoms, vaginal atrophy (Klingsberg S. A. et al., Int J Women's Health, 1: 105-111 (2009); Mac Bride, M. B., Mayo Clin. Proc., 85:87-94 (2010); Smith, A. L et al., Discov. Med., 10:500-510 (2010)); colon cancer (Oduwole O. O. et al., J. Steroid. Biochem. Mol. Biol. 87:133-140 (2003)), neuronal diseases (Behl, C. et al., Biochem. Biophys. Res. Commun., 216:473-482 (1995)) like Alzheimer (Pike, C. J., Front Neuroendocrinol., 30:239-258 (2009)) and Parkinson's disease (Bourque, M. et al., Front Neuroendocrinol., 30:142-157 (2009)), depression (Schmidt, P. J. et al., Ann. N.Y. Acad. Sci., 179:70-85 (2009)) anxiety, hypercholesterolemia (Karjalainen, A. et al., Arterioscler. Thromb. Vasc. Biol., 4:1101-1106 (2000)), cardiovascular diseases (Traish, A. M. et al., J. Androl., 30:477-494 (2009); Xing D. et al., Arteriosder. Thromb. Vasc. Biol., 29:289-295 (2009)); hair loss (Mooradian, A. D. et al., Endocr. Rev., 8:1-28 (1987); Georgala S. et al., Dermatology 208:178-9 (2004)), non-insulin-dependent diabetes mellitus (Ferrara A. et al., Diabetes Care 25: 1144-1150 (2001)), rheumatic diseases and inflammatory diseases (Cutolo, M. et al., Ann. N.Y. Sci., 1193:36-42 (2010); Islander U. et al., Mol. Cell. Endocrinol. Doi:10.1016/j.mce. 2010.05.018 (2010)). More recently, non-steroidal 17β-HSD2 inhibitors from different compound classes have been described (Wetzel M. et al., Bioorg. Med. Chem., 19: 807-815 (2011); Wetzel M. et al., J. Med. Chem., 54: 7547-7557 (2011). Wetzel M. et al., Eur. J. Med. Chem., 47:1-17 (2012). Xu K. et al., Eur. J. Med. Chem., 46: 5978-5990 (2011); Xu K. et al., Letters in Drug Design & Discovery, 8: 406-421 (2011); Al-Soud Y. A. et al, Arch. Pharm. (Weinheim), 345: 610-621 (2012); Marchais-Oberwinkler et al. S., J. Med. Chem., 56:167-181 (2013); Perspicace E. et al., Eur. J. Med. Chem., 69:201-215 (2013); Perspicace E. et al., Molecules, 18: 4487-4509 (2013)). Further, selective 17bHSD1 inhibitors with N-benzyl-N-methyl(phenyl)thiophene-carboxamide, N-benzyl-N-methyl-(phenyl)-1,3-thiazole-carboxamide and N-benzyl-N-methyl-biphenyl-3-carboxamide structure are known from WO2012117097.\nRelated Structures: Thakar K. A. and Padhye A. M., J. Ind. Chem. Soc. LXI: 715-716 (1984) discloses certain mono- and dihalogeno-hydroxyphenyl-furyl-ketones, and mono- and di-halogeno-hydroxyphenyl-thienyl-ketones in the synthesis of furyland thienyl-substituted benzisoxazoles.\nTakayama T. et al., Bioorg. & Med. Chem. Lett. 20: 108-111 (2010) and WO2005/123671 disclose (3-amino-4-carboxamido-5-(4′-phenoxy)phenyl-pyrrol-2-yl)-phenyl-ketone derivatives that are inhibitors of lymphocyte-specific kinase and are useful as immunsuppresive agents, e.g. for the treatment of inflammatory diseases and of organ transplant rejection.\nEP0801058A1 and U.S. Pat. No. 5,817,691 disclose arylthio-, arylsulfinyl- and arylsulfonylpyrrole compounds having insecticidal activity.\nWO2011/079772 discloses the one-pot synthesis of certain 2,5-disubstituted thienes."}
{"text": "The background of the present disclosure and the illustrative embodiments disclosed herein are described in the context of identifying known audio recordings encountered during an outbound telephone call, for example during a call placed from a contact center. However, the present invention has applicability to the identification of any segment of audio or an image (as used herein, the term “image” is intended to encompass both still and moving images), regardless of the type or source of the audio or image, and regardless of in what circumstances the audio or image is encountered. Furthermore, the present invention also has applicability to the identification of any segment of data such as, for example, data obtained from any type of sensor. Therefore, as used herein, the term “dataset” shall encompass a collection of any type of data, whether comprising audio, image, or other type of data.\nIn a classic contact center scenario, outbound calls are made either automatically (by a class of devices known as “automated dialers” or “autodialers”) or manually. A number of human “agents” are available to join into calls that are determined to reach a live person at the called end. In this way, efficiencies are obtained by not having agents involved in a call until it is determined that there is a live person at the called end with whom the agent may speak. The use of automated equipment to monitor the telephone line during the outbound call is referred to as call progress analysis (CPA). CPA is a class of algorithms that operate on audio and network signaling during call setup. The goal of CPA is to determine the nature of the callee, or the outcome of call setup to an external network (traditional public switched telephone network or Voice over Internet Protocol (VoIP)). Specifically, when a call or session is being established, the caller or initiator must determine whether it was answered by a live speaker, if the line is busy, etc. When the caller is an automated application, such as an automated dialer or message broadcasting system, CPA algorithms are used to perform the classification automatically. CPA is used to interpret so-called call-progress tones, such as ring back and busy, that are delivered by the telephone network to the calling entity. Traditional CPA is performed using low- and high-pass frequency discriminators together with energy measurements over time to qualify in-band signaling tones.\nAnother method for classifying audio on an outbound call is known as Voice Activity Detection (VAD), which is a class of audio processing algorithms that identify where speech is present in an audio stream. The detected speech may originate from any source, including a live speaker or a prerecorded message. Modern VAD algorithms use spectral analysis to distinguish the utterance of a primary speaker from background noise.\nA subclass of CPA algorithms that extract speaking patterns using VAD, and determine whether the patterns originate from a live speaker or a prerecorded message, is known as Answering Machine Detection (AMD). By identifying calls that do not connect to a live speaker, an accurate AMD algorithm can significantly increase throughput of an automated dialer. However, false positives from AMD lead to silent or abandoned calls, causing revenue loss for the contact center, and negative impressions amongst the public. The quality of an AMD algorithm is a function of the accuracy and response time, and some regions of the world (notably the U.S. and U.K.) impose strict legal requirements on both.\nAMD is not an exact science, and the optimal approach is an open problem. To achieve acceptable accuracy, speed, and flexibility, AMD algorithms use a combination of heuristics and statistical models such as neural networks to classify an utterance as live or pre-recorded. Although many commercial AMD systems available on the market report high accuracy rates in the marketing literature (e.g., 95% or more), there is no independent auditor for these figures, and the actual accuracy rate is typically much lower in practice (e.g., 80% or less), as reflected by continued widespread complaints. A general ban has been proposed by some consumer advocacy groups, and some contact centers simply cannot use AMD because of its limitations.\nA relatively new science of audio identification is known as Acoustic Fingerprinting, in which a system generates a “fingerprint” of a candidate audio stream, and compares it against a database of known fingerprints, analogous to human fingerprinting used in forensics. In this context, a “fingerprint” is a condensed digest of an audio stream that can quickly establish perceptual equality with other audio streams. A database of known fingerprints may associate known fingerprints with meta-data such as “title”, “artist”, etc. The past ten years have seen a rapidly growing scientific and industrial interest in fingerprinting technology for audio and images. Applications include identifying songs and advertisements, media library management, and copyright compliance.\nVarious acoustic fingerprinting algorithm classes have been proposed, and the most prevalent today are those based on either “landmarks” or “bitmaps”. Landmark-based algorithms extract discrete features from an audio stream called “landmarks”, such as spectral peaks, sudden changes in tone, pitch, loudness, etc. The optimal choice of landmark is an open question guided mostly by heuristics. The acoustic fingerprint is stored as a sequence of data structures that describe each landmark. At runtime, landmarks extracted from a candidate audio stream are compared to a database of fingerprints based on a distance metric.\nBitmap-based algorithms analyze an audio stream as a sequence of frames, and use a filter bank to quantize each frame into a bit vector of size N, where N is typically chosen for convenience as the number of bits in a C-style integer, e.g. N∈{8, 16, 32, or 64}. A popular and well-studied example is known as the “Haitsma-Kalker algorithm”, which computes a binary bitmap using a filter that compares short-term differences in both time and frequency. The Haitsma-Kalker Algorithm has been well-studied in the literature. It's inventors, Jaap Haitsma and Ton Kalker, have published a report of use of the Haitsma-Kalker Algorithm and the comparison of binary acoustic fingerprint bitmaps to identify three (3) second recordings of songs from a database of millions of songs (Haitsma and Kalker, “A Highly Robust Audio Fingerprinting System,” Journal of New Music Research, Vol. 32, No. 2 (2003), pp. 211-221). The complete acoustic fingerprint is stored as a sequence of bit vectors, or a bitmap. As illustrated in FIG. 1A-C, there are shown three images of an audio stream containing a message from a telephone network saying “This number has been disconnected”. FIG. 1A shows the original audio wave signal, with 1.5 seconds of audio sampled at 8000 KHz. FIG. 1B shows a spectrogram of the original audio input signal, with dark regions indicating high energy at a particular frequency. FIG. 1C shows a binary acoustic fingerprint bitmap created using the Haitsma-Kalker algorithm, with height N=16. The height is determined by the number of bits computed at each frame, and the width is determined by the number of frames in the audio stream. At runtime, the bitmap computed from a candidate audio stream is compared to a database of bitmaps based on the number of non-matching bits, also known as the Hamming distance.\nThe use of bitmap matching and the process of acoustic fingerprinting is a powerful emerging tool in the science of audio recognition; however, it is computationally intense and requires several seconds of sampled audio to make a match in many cases. This delay makes it not well suited for use in call progress analysis. Accordingly, there remains a need for faster and more accurate systems and methods for identifying audio, both in the general case and during an outbound call attempt."}
{"text": "1. Field of the Invention\nThe present invention relates to a photonic crystal fiber (PCF) preform and a photonic crystal fiber manufactured using the same.\n2. Description of the Related Art\nA photonic crystal fiber (PCF) is a special kind of optic fiber. A general single-mode optical fiber is made up of a core material which consists of glass and germanium or phosphorus. On the other hand, the photonic crystal fiber is made up of a single, solid and substantially transparent material 1, such as silica glass, and periodically spaced air holes 2 running along the fiber, as shown in FIG. 1. Difference in dielectric constant between the silica glass material 1 and the periodically spaced air holes 2 forms a photonic band gap, which serves to prevent movement of a ray in a specific direction as in an electronic band gap of a semiconductor. In other words, the ray passes through the photonic band gap only when the ray satisfies conditions under which it may pass through.\nThe photonic crystal fiber has many important technical properties. For example, the photonic crystal fiber may support a single mode over a wide range of wavelengths, and may provide a large mode region. Consequently, the photonic crystal fiber is capable of transmitting high optical power and providing large phase dispersion in a telecommunication wavelength of 1.55 μm. Furthermore, the photonic crystal fiber is increasingly used as a device for controlling nonlinearity or polarization. It is expected based on recent reports on these functions of photonic crystal fiber, that photonic crystal fiber will be widely applied to optical communication and optical-related industries in the near future.\nA photonic crystal fiber must originate from a photonic crystal fiber perform, which is drawn into a photonic crystal fiber having an efficiently long length and retaining its\nGenerally, the photonic crystal fiber preform may have structures schematically shown in FIGS. 2a and 2b. \nThe photonic crystal fiber preform of FIG. 2a is made up of a circular silica glass rod 10 and a plurality of circular air holes 11 longitudinally formed through the circular silica glass rod 10 in a photonic lattice structure. Index of refraction of the photonic crystal fiber preform with the above-stated construction is distributed as illustrated in FIG. 2b. \nThe conventional photonic lattice structures, however, limit photonic crystal fiber preform design. Furthermore, the air holes may be deformed when being drawn into desired photonic crystal fibers. The photonic crystal fibers are, as a result, not manufactured according to the originally intended design. Also, the material constituting a core of the photonic crystal fiber is limited as to its properties."}
{"text": "Lithographic projection apparatus (tools) can be used, for example, in the manufacture of integrated circuits (ICs). In such a case, the mask contains a circuit pattern corresponding to an individual layer of the IC, and this pattern can be imaged onto a target portion (e.g. comprising one or more dies) on a substrate (for example, but not limited to a silicon wafer) that has been coated with a layer of radiation-sensitive material (resist). In general, a single wafer will contain a whole array of adjacent target portions that are successively irradiated via the projection system, one at a time. In one type of lithographic projection apparatus, each target portion is irradiated by exposing the entire reticle pattern onto the target portion in one go; such an apparatus is commonly referred to as a wafer stepper. In an alternative apparatus—commonly referred to as a step-and-scan apparatus—each target portion is irradiated by progressively scanning the mask pattern under the projection beam in a given reference direction (the “scanning” direction) while synchronously scanning the substrate table parallel or anti-parallel to this direction; since, in general, the projection system will have a magnification factor M (generally <1), the speed V at which the substrate table is scanned will be a factor M times that at which the mask table is scanned. More information with regard to lithographic apparatus as here described can be gleaned, for example, from U.S. Pat. No. 6,046,792, incorporated herein by reference.\nIn a manufacturing process using a lithographic projection apparatus, a mask pattern is imaged onto a substrate that is at least partially covered by a layer of radiation-sensitive material (resist). Prior to this imaging step, the substrate may undergo various procedures, such as priming, resist coating and a soft bake. After exposure, the substrate may be subjected to other procedures, such as a post-exposure bake (PEB), development, a hard bake and measurement/inspection of the imaged features. This array of procedures is used as a basis to pattern an individual layer of a device, e.g. an IC. Such a patterned layer may then undergo various processes such as etching, ion-implantation (doping), metallization, oxidation, chemo-mechanical polishing, etc., all intended to finish off an individual layer. If several layers are required, then the whole procedure, or a variant thereof, will have to be repeated for each new layer. Eventually, an array of devices will be present on the substrate (wafer). These devices are then separated from one another by a technique such as dicing or sawing. Thereafter, the individual devices can be mounted on a carrier, connected to pins, etc. Further information regarding such processes can be obtained, for example, from the book “Microchip Fabrication: A Practical Guide to Semiconductor Processing”, Third Edition, by Peter van Zant, McGraw Hill Publishing Co., 1997, ISBN 0-07-067250-4, incorporated herein by reference.\nThe lithographic tool may be of a type having two or more substrate tables (and/or two or more mask tables). In such “multiple stage” devices the additional tables may be used in parallel, or preparatory steps may be carried out on one or more tables while one or more other tables are being used for exposures. Twin stage lithographic tools are described, for example, in U.S. Pat. No. 5,969,441 and WO 98/40791, incorporated herein by reference.\nThe photolithography masks referred to above comprise geometric patterns corresponding to the circuit components to be integrated onto a silicon wafer. The patterns used to create such masks are generated utilizing CAD (computer-aided design) programs, this process often being referred to as EDA (electronic design automation). Most CAD programs follow a set of predetermined design rules in order to create functional masks. These rules are set by processing and design limitations. For example, design rules define the space tolerance between circuit devices (such as gates, capacitors, etc.) or interconnect lines, so as to ensure that the circuit devices or lines do not interact with one another in an undesirable way.\nOf course, one of the goals in integrated circuit fabrication is to faithfully reproduce the original circuit design on the wafer (via the mask). Another goal is to use as much of the semiconductor wafer real estate as possible. As the size of an integrated circuit is reduced and its density increases, however, the CD (critical dimension) of its corresponding mask pattern approaches the resolution limit of the optical exposure tool. The resolution for an exposure tool is defined as the minimum feature that the exposure tool can repeatedly expose on the wafer. The resolution value of present exposure equipment often constrains the CD for many advanced IC circuit designs.\nFurthermore, the constant improvements in microprocessor speed, memory packing density and low power consumption for micro-electronic components are directly related to the ability of lithography techniques to transfer and form patterns onto the various layers of a semiconductor device. The current state of the art requires patterning of CD's well below the available light source wavelengths. For instance the current production wavelength of 248 nm is being pushed towards patterning of CD's smaller than 100 nm. This industry trend will continue and possibly accelerate in the next 5-10 years, as described in the International Technology Roadmap for Semiconductors (ITRS 2000).\nThis continued demand for improved performance has resulted in the development of various techniques aimed at improving resolution. Such techniques are typically referred to as Resolution Enhancement Techniques (RET's) and comprise a very wide range of applications. Examples include: light source modifications (e.g. Off-Axis Illumination), use of special masks, which exploit light interference phenomena (e.g. Attenuated Phase Shift Masks, Alternating Phase Shift Masks, Chromeless Masks, etc.), and mask layout modifications (e.g. Optical Proximity Corrections).\nOf the foregoing techniques, dipole illumination is one of the most attractive RET candidates due to its high image contrast for dense pitches and superior resolution capabilities. As is known, dipole illumination is an extreme case of OAI and is capable of providing enhanced imaging contrast with improved process latitude for very low K1 imaging.\nHowever, one of the limitations associated with dipole illumination is that a single illumination only enhances resolution for features that are orthogonal to the illumination pole axis. As a result, in order to take full advantage of dipole illumination during wafer printing, the mask pattern must be decomposed into horizontal and vertical orientations. Once the mask pattern is converted in this manner, a Y-pole exposure is utilized to image the horizontally oriented features, and a X-pole exposure is utilized to image the vertically oriented features. One important aspect of dipole illumination is that when imaging the horizontally oriented features, the vertically oriented features must be protected (i.e., shielded) so the vertically oriented features are not degraded. The opposite is true when vertically oriented features are imaged (i.e., the horizontally oriented features must be protected).\nFIG. 1 illustrates the basic concepts of double dipole imaging. As stated, typically there are at least two exposures when utilizing dipole illumination. In the first exposure, the X dipole aperture 10 provides a maximum aerial image intensity (i.e., maximum modulation) for the vertical portion of the lines 12 to be printed. The resulting image profile is illustrated by line 24 in FIG. 1. In the second exposure, which utilizes the Y-dipole aperture 16, there is no image modulation for lines 12. It is noted, however, that during the second exposure using the Y-dipole aperture, the vertical portions of the lines 12 need to be shielded so that the vertical features formed during the first exposure are not degraded during the second exposure. FIG. 1 illustrates shielding the lines 12 with shields 15, each of which is 20 nm wide in the horizontal direction. As a result, when exposing the horizontal lines using the Y dipole aperture, there is substantially no imaging (i.e., modulation) of the vertical features 12. The aerial image is a DC modulation as shown by line 17 in FIG. 1, which corresponds to the 20 nm shielding. The final aerial image intensity, which is represented by line 14 in FIG. 1, corresponds to the sum of the first exposure using the X dipole aperture and the second exposure using the Y dipole aperture.\nIt is further noted that, assuming the exposure energy is constant, increasing the width of the shielding from a 20 nm shield 15 to a 40 nm shield 20 for the vertical lines 12 causes the minimal intensity level of the resulting image to shift to a lower level. This is represented by line 22 in FIG. 1, which represents the aerial image associated with the vertical portions of the features. As shown, the aerial image 22 is just a DC modulation. However, it is lower than the DC modulation 17 associated with the 20 nm shield. As a result, the composite image 19 formed utilizing the 40 nm shielding provides better imaging results than the composite image 14 formed utilizing the 20 nm shielding.\nAs a result of the need to separate the horizontally and vertically oriented features, one of the challenges for the lithographer, when utilizing dipole illumination, is determining how to convert the original IC design data into its horizontal or vertical pattern components and generate two masks for the dual exposure process that can take full advantage of the dipole imaging performance. One factor that reduces performance and which should be considered when generating the mask patterns is background light due to lens flare or scattering. As is known, lens flare results in unwanted background light (i.e., noise) that degrades the image contrast at the image plane. Thus, it is desirable to reduce “flare” as much as possible. This is especially true when utilizing dipole illumination techniques due to the multiple exposures associated therewith.\nThe “aerial image with flare” is equal to the “aerial image without flare” convolved with a point-spread function (PSF) plus the scattering. The foregoing can be expressed as:Iflare(x,y)=InoflarePSFflare+Inoflare(I−TIS)  (1)\nwhere TIS is the total integrated scattering (TIS) for lens having a surface roughness with a Gaussian-like distribution. Under such conditions, TIS can be expressed as:TIS=[(4πσ cos θ)/λ]2  (2)\nwhere λ is the wavelength of the exposure tool, σ is the rms roughness of the lens, and θ is the scattering angle. As a result of current lens making capabilities, which result in lens exhibiting extremely low surface roughness, the foregoing equation can be approximated as:TIS˜1/λ2  (3)\nEquation (3) makes clear that as the wavelength of the exposure tool is reduced, the amount of scattered light increases significantly. For example, the total integrated scattering (TIS) of light for an exposure tool having a wavelength of 193 nm is approximately 1.65 times greater that the TIS associated with an exposure tool having a wavelength of 248 nm.\nIt is noted that the first term is equation (1) is the “diffuse halo” which causes the focused image to spread out. The second term in equation (1) is the contribution due to scattering. The overall effect is an unwanted DC background light that reduces the aerial image contrast. Furthermore, besides the negative impact on image contrast, flare is also unevenly distributed across the scanning slit and is not uniform with the exposure field, which can cause intrafield CD variations. Therefore, protecting features and reducing background stray light becomes increasingly critical. The issue of how to reduce or negate the effects of background stray light becomes even more important as the wavelengths of the exposure tools are reduced.\nCurrently, one known technique for reducing the negative effects of flare comprises the step of adding solid chrome shielding on the large areas of the mask pattern (i.e., background portions) that do not contain any geometry (i.e., features). As shown in FIGS. 2a and 2b, when utilizing dipole illumination, the solid chrome shielding, referred to as background light shielding (BLS), is applied to the background areas in both the horizontal mask and the vertical mask. The solid chrome shield functions to protect the background area during both exposures. FIG. 2a illustrates an example of the use of this shielding technique in conjunction with the printing of horizontally oriented features 29 utilizing the Y dipole 16. As shown in FIG. 2a, each of the vertical features 27 are provided with shielding 210 (i.e., main feature shielding) in the manner discussed above in conjunction with FIG. 1. In addition, the solid chrome shield 220 is provided in the background area where there are no features to be imaged on the wafer. In a similar manner, FIG. 2b illustrates the vertical mask, in which the horizontal oriented features are shielded, while the vertical features are printed. As shown, the vertical mask also includes a solid chrome shield 220 disposed in the background area. It is further noted that both the horizontal mask and the vertical mask contain assist features 103 (e.g., scatter bars).\nHowever, as a result of such background shielding 220, when utilizing a positive resist, the intensity in the background areas becomes too low to completely clear the resist. FIGS. 3a and 3b illustrate a simulated resist pattern corresponding to the portion of the masks of FIGS. 2a and 2b defined by area 30, which includes the solid chrome shielding 220. The simulation was performed assuming NA (numerical aperture)=0.75, ArF double exposure x-pole, y-pole, σouter/σinner=0.89/0.65. As is shown by FIGS. 3a and 3b, portions of the resist 221 in the background areas remain after illumination using the vertical and horizontal masks. As a result, a third exposure utilizing a trim mask is necessary in order to completely remove the resist from the background shielded areas. Thus, such a solution for reducing the effects of flare is undesirable as it results in an increase in the number of exposures and masks required for imaging the wafer. Referring to FIG. 3a, the areas indicated by reference numeral 51 correspond to the areas where resist remains after the double exposures, and these areas are contrasted against areas of either the vertical or horizontal mask that had chrome disposed thereon (i.e., either feature or shielding).\nFurthermore, the foregoing solid chrome shielding technique can also negatively interfere with assist features, such as scatter bars, and cause the assist features to print underneath the shielding of either the horizontal or vertical mask, as also illustrated in FIGS. 3a and 3b. For example, referring to FIG. 3b, as shown in the resist simulation, the assist features 103, which are intended to be sub-resolution, are printed as a result of the BLS 220. This problem imposes an additional constraint on the placement of assist features, which can prevent the assist features from being placed in the optimal position, thereby causing a reduction in printing performance.\nAccordingly, there exists a need for a method for negating the effects of flare in the exposure process which does not result in an increase in the number of exposures and masks required for imaging the wafer, and which does not impact the use and/or placement of assist features in the mask."}
{"text": "The embodiments described herein relate to a coil receiving MR signals from a subject, and an MRI (magnetic resonance imaging) system including the coil.\nA reception coil for receiving MR signals from a subject is employed in MRI apparatus.\nThe reception coil known from, for example Japanese Patent Application No. Hei 08-252234, has a problem in which application of low frequency noise to the reception coil results in artifacts such as ghosts appearing in an MR image."}
{"text": "The invention relates to a height adjustment device of a vehicle seat, which comprises a control valve means serving for controlling the degree of filling of at least one compressed air actuator holding the seat body at a desired height and thereby providing a vertical elastic supporting action, such control valve means being so provided with at least one actuating member for cooperation with at least one control cam slide that in accordance with use dependent alteration of the weight to be borne by the at least one compressed air actuator a relative switching over movement causing the switching over of the control valve means takes place between a component bearing the control cam slide and another component bearing the associated actuating member.\nThe invention furthermore relates to a vehicle seat fitted with such a height adjustment device, comprising a seat body, which is arranged on a seat bracket fitted with a height adjustment device, such seat bracket having at least one compressed air actuator, with the aid of which the seat body may be set at different heights and is able to be vertically resiliently supported at the respectively set height, the seat bracket having at least two support elements changing the distance between them during vertical movement of the seat body, one support element of the seat bracket being kinematically coupled with at least one control cam slide, which so cooperates with at least one actuating member kinematically coupled with the other support element and associated with a control valve means serving for control of the degree of filling of the at least one compressed air actuator that in the case of an alteration in the use dependent weight carried by the at least one compressed air actuator a relative switch over between the at least one control cam slide and the components bearing the actuating member associated with the control cam slide and the actuating member associated with it.\nSuch a height adjustment device with an associated vehicle seat is disclosed in the German patent publication DP 4335199. The seat bracket of the vehicle seat in this case comprises two support elements, coupled together to form a scissor-like structure, which alter their angular position when the height of the seat body changes, for example owing to a varying weight load. The change in angle is detected by two control valves, whose actuating members control cam slide down two curved control cam slides while performing an arcuate switching over movement and, if they are deflected sufficiently, cause compressed air to be supplied to the compressed air actuator or let off therefrom. Accordingly it is possible to ensure that the seat height is automatically set to the right height independently of weight.\nOwing to their quite complex structure the known height adjustment devices appear to be in need of improvement. This also applies for alternative designs of height adjustment devices and the vehicle seats fitted with them, which are described in the German patent publications DE 19705010 A1, DE 3517503 C2 and DE 3517505 C2."}
{"text": "This invention relates to a synchronization signal processing system for use in a mobile communication network which comprises a plurality of mobile service switching centers and a plurality of base transceiver stations and is operable in a time division fashion.\nThe mobile communication network has an overall service area which is divided into cells or radio zones assigned with the base transceiver stations, respectively, and in which a plurality of mobile stations are present, namely, either moving or staying standstill, at a time. Each mobile station may be either a portable telephone device carried by a user or a subscriber's terminal installed in an automobile or in a like mobile vehicle and is movable from a first zone of the cells to a second zone of the cells.\nIt is possible to understand that each mobile service switching center is connected to a plurality of fixed subscriber substations either directly or through at least one exchange office. Some of the mobile service switching centers are connected to the base transceiver stations. More particularly, each of such mobile service switching centers is connected to a certain number of base transceiver stations.\nThe mobile service switching centers are connected to one another by wired communication lines. The mobile service switching centers and the base transceiver stations may be connected through wired communication lines. Among the overall service area, some of the cells are often referred to collectively as a radio communication area when assigned to the base transceiver stations which are served by one of the mobile service switching centers.\nEach base transceiver station is for transmitting and receiving radio message signals to and from at least one of the mobile stations that is currently present in the cell assigned with the base transceiver station under consideration. For use in time division multiple access (TDMA), the radio message signals are carried by a radio carrier signal of a radio frequency in a plurality of time slots. A predetermined number of such time slots are successively arranged in a frame in the manner known in the art.\nWhen a particular station of the mobile stations moves between the first and the second zones assigned with first and second stations of the base transceiver stations, the first and the second stations use different radio frequencies and different time slots in transmitting and receiving the radio message signals to and from the particular station. The first and the second stations may be connected either to one or to two of the base transceiver stations. In either event, the particular station is inevitably subjected to a handover processing between the first and the second stations. It is therefore desirable to preliminarily synchronize the frames and the time slots in the base transceiver stations in order to reduce a time necessary for such a handover processing as a handover processing time.\nIn the manner which will later be described, a conventional synchronization signal processing system comprises an individual synchronization signal generating circuit in each mobile service switching center. When connected to such a mobile service switching center, the base transceiver station can generate synchronized frames and synchronized time slots for the mobile stations which are currently present in the radio communication area served by the base transceiver station under consideration.\nA little more in detail, the synchronization signal generating circuit comprises first and second time division switches, each comprising controllable connection paths and producing a switch trouble signal when a trouble occurs therein. A controller device is cross connected to the first and the second time division switches and is supplied with the switch trouble signal to control the connection paths of one of the first and the second time division switches that is not producing the switch trouble signal and serves as an active switch with the other of the first and the second time division switches used as a standby switch. A synchronization signal generator is connected to the active switch to supply a synchronization signal to the connection paths of the active switch. Output trunk circuits are connected to the connection paths of the first and the second time division switches to supply the synchronization signal to at least one of the output trunk circuit from the connection paths controlled by the controller device to the base transceiver stations served by mobile service switching center in question.\nIt is liable that the synchronization signal generator is involved into a trouble. First and second synchronization signal generators are therefore cross connected to the first and the second time division switches. Alternatively, it is possible to understand that the first and the second synchronization signal generators are connected to the active switch. In either event, each synchronization signal generator produces a generator trouble signal when a trouble occurs therein. Supplied with the generator trouble signal, the controller device controls the connection paths of the active switch to supply the output trunk circuits with the synchronization signal generated by one of the first and the second synchronization signal generators that is not producing the generator trouble signal.\nAs a consequence, the conventional synchronization signal processing system can deal with troubles that may occur in the time division switches and/or in the synchronization signal generators. It is, however, impossible to keep the phase of the synchronization signal when the first and the second synchronization signal generators are switched from one to the other."}
{"text": "It is described in German Published Patent Application No. 35 36 820 and ISO 14819 that traffic messages in the form of digitally encoded messages plus radio programs may be broadcast over radio frequencies to describe traffic-relevant situations, in particular traffic disturbances in the highway system. These TMC (Traffic Message Channel) traffic messages include location information about the location of a traffic disturbance in an encoded form.\nISO 14819 also describes so-called multisequence messages in which traffic information is transmitted in several groups of the RDS signal, but the several groups which include traffic information must always be transmitted in direct succession.\nGerman Published Patent Application No. 199 05 893 describes an expansion of traditional traffic messages according to the TMC standard. It is provided there that a supplementary location description, which is announced in a header preceding the actual message, is to be added to the standardized messages, which regularly contain a location code and thus a reference to a location of a traffic-relevant event. Thus, a location description is no longer limited merely to highway junctions, highway intersections and interchanges and the like and/or the sections in between that are encoded in the TMC location database but instead it allows a further description of the event location.\nA more accurate localization of an event location is also the subject of German Published Patent Application No. 100 15 935. It is proposed there that in addition to a section of road affected by a traffic disturbance, which may be defined by an adjacent location encoded in the TMC location database, a portion of a section or comparable linear parameters may also be transmitted, permitting a more accurate localization of the event location on the encoded section of road.\nThe traffic situations to be transmitted via a traffic message may be simple or complex; for example, “10 km backed-up traffic” is a simple description of the situation and “10 km backed-up traffic, construction site, lane closure, average speed=20 km/h” is a complex description of the situation. Such complex situation descriptions may be described by so-called “multisequence messages,” i.e., multiple indexed successive individual messages in TMC (Traffic Message Channel, as specified in ISO 14819).\nOne disadvantage here is that all the individual events of a complex situation description must always be based on the same location, i.e., the same section of road, so that the message may be displayed on the terminal as a complex situation description. Furthermore, all the events of a situation must be sent at the same point in time. Although updating is possible, all messages, including all the events they contain, must always be updated. It is impossible to append additional events to a message already sent. Expanded complex example: “Between Laatzen junction and Hildesheim junction 10 km backed-up traffic, construction site, lane closure, average speed=20 km/h.” The individual events of “backed-up traffic,” “construction site,” “lane closure,” and “average speed=20 km/h” must be based on the same location, namely in this case the same section of road between junctions. If the events overlap or if they are based on different neighboring locations, multiple separate messages must be transmitted. It is very complex to combine the messages at the terminal end to allow a compact presentation.\nIt is also a disadvantage that it is impossible to correlate messages originating from different sources. In the future it may be expected that situation descriptions of differing content will be supplied by different sources. For example, traffic disturbances such as congestion or accidents are compiled by the police via the state reporting offices, while long-term status information such as construction sites or gridlock is supplied by third-party providers—possibly even as a paid service. Example: real situation: “10 km congestion and 5 km construction site.” It is assumed that the construction site will remain in existence for a longer period of time and the message will be transmitted regularly by provider X, e.g., a radio station. Congestion occurs spontaneously and is reported by a state reporting office for a relatively short period of time. With the digitally encoded TMC traffic messages currently being transmitted by radio, there is no possibility of connecting two individual messages to form one complex message."}
{"text": "Arc discharge in aqueous electrolytes (for example, welding under seawater), is widely used in engineering and construction, and is at present the only known form of stationary plasma discharge in liquid media. In recent years, such discharge was also used in different physicochemical studies and in the synthesis of various materials. The specific feature of arc discharge in liquid media is the localization of a plasma region near the electrode ends and a “falling” form of volt-ampere characteristic as illustrated in FIG. 1.\nIn a gaseous phase, different kinds of discharges can be implemented, the external manifestation and electrical parameters of which are connected with a wide range of technical characteristics for devices used in their implementation and a variety of elementary processes determining the conditions of current passage through gas. The essential feature of electric discharge development in the gaseous phase is a profound effect of the properties of the gas medium on the current passages through the gas.\nUnder usual conditions, the concentration of charge carriers (electrons and ions) in the gas is very low: a gas is a very good dielectric. For a gas to have a high electrical conductivity (as a result of ionization) it is necessary for a high quantity of charge carriers to be present, requiring in turn a great quantity of energy. Gases have a steady electric conductivity when there is equilibrium between the origination and disappearance of charges. Thus, to create a means by which high electrical conductivity in a gas can be achieved through substantially lower energy requirements than has been taught in the prior art is highly desirable.\nIf the rate of movement of electrical charges is proportional to the field strength, the conductivity of gas approximately obeys Ohm's law (FIG. 2, section a). With increasing field strength, the decrease of electrical charges begins to have an influence (FIG. 2, section b) because of the migration of the charges to the electrodes. Further increases of the electrical field strength result in a steep increase of current due to the start of collision ionization (FIG. 2, section c). In spite of the avalanche-like character of current increases, the existence of external ionizer(s) is needed to sustain the electrical discharge, and the discharge remains being as not self-sustained (region 1). Eventually, a point is reached where for each electron leaving the cathode, one or more electrons arrive at the anode, in a phenomenon known as breakdown discharge (glow discharge or plasma discharge). This causes a self-sustained electrical current from the cathode to the anode. However, the current state-of-the-art process requires a large amount of energy to reach this self-sustaining threshold. Since high energy requirements directly and indirectly decrease the overall economy of the model, the requirement of high energy is undesirable. Therefore, it is highly desirable to have a new process having low energy demand in which the transition from non self-sustained discharge to self-sustained discharge (glow discharge) would occur with a low-energy input.\nAgain referring to FIG. 1, which illustrates the prior art, the voltage-current characteristic curve for glow discharge preferably comprises three sections, referred to for the sake of clarity as subnormal section or subnormal mode (FIG. 2, section d), normal section or normal mode (FIG. 2, section e) and abnormal section or abnormal mode (FIG. 2, section f).\nFurther increase of current density on the cathode causes the appearance of electric arc, as well as a drastic change of the main characteristics of the discharge (FIG. 2, section g).\nIt should be noted that the appearance or threshold of discharges in the gas phase depends considerably on the pressure of the gas. Thus, in the case of a uniform field of breakdown voltage (self-maintained discharge initiation voltage) the threshold is determined by the product of pressure by the distance between the electrodes, according to Paschen's Law. Pachen determined that breakdown voltage is determined by the following equation:\n  V  =            a      ⁡              (        pd        )                            ln        ⁡                  (          pd          )                    +      b      where V is the breakdown voltage in Volts, p is the pressure in atmospheres, d is the gap distance in meters, and a and b are constants that depend upon the particular gas between the electrodes. Thus, in contrast to liquids, which are relatively incompressible, different forms of electric discharge can be implemented in gases by varying the pressure of the gas between the electrodes.\nMoreover, when ultrasonic cavitation, a sort of “cold boiling” resulting from the creation and collapse of zillions of microscopic bubbles in the liquid caused by ultrasonic waves, is implemented within a liquid, its phase composition and physical properties abruptly change, which can lead to some specific features for the formation of electric discharges within the liquid. In the region of intense cavitation, a gaseous component is formed which represents a significant fraction of the liquid. Therefore it can be assumed that the conditions for electric breakdown into the cavitation region should become easier, and the initiation of different forms of discharge could start through use of this invention. By varying the parameters of an ultrasonic field, it is possible to influence the processes of plasma glow within a cavitating liquid.\nThe prior art has several examples of attempts to resolve this problem.\nHowever, few patent applications or patents work in the abnormal mode. In abnormal mode, also known as abnormal glow, effectively all of the gas molecules must be ionized to provide charge carriers for the current. Typically, the gas molecules are ionized multiple times meaning that more than one electron has been freed for most of the gas molecules. This creates a relatively uniformly distributed plasma across the electrodes. A higher density (or pressure) of gas molecules, on the other hand, would lead to a normal mode, or normal glow discharge. In this region, fewer than all of the molecules are ionized. This creates a situation where plasma forms in a relatively small region between the electrodes. A plasma discharge of this type can lead to concentrated energy in a relatively small area and possibly lead to electrode damage. Therefore, it is preferable to work in the abnormal mode.\nThose patent applications or patents that do work in the abnormal mode, like U.S. Pat. No. 5,068,002, to Monroe, do not use an electrode as the radiator, in the same way that the instant application uses it, whereby the current application discloses a very low energy consumption jointly with a very low voltage to initiate and maintains a volumetric discharge which generates operational advantages in term of achieving the goals of this application. Monroe describes an ultrasonic glow discharge surface cleaning apparatus for abrading contaminants from the surface of a work piece using plasma glow discharge.\nFor example, in US Patent Application 2004/0265137 A1 to Bar-Gadda, a method is proposed for hydrogen production from water or steam by means of plasma discharge excited in the UHF, radio- or low-frequency range, as well as with arc discharge. This application describes the injection of water molecules into plasma discharge.\nU.S. Pat. No. 7,070,634 B1 A1 to Wang describes a plasma apparatus for converting a gaseous mixture of water vapor and hydrocarbons into hydrogen.\nUS Patent Application 2006/0060464 to Chang teaches a fluid phase contained in a reactor, within which electrodes (anode and cathode) are placed. A flow of gas bubbles is introduced or generated in the medium in the region adjacent to the cathode. The potential difference necessary for the initiation of glow discharge and for the ionization of gas molecules in the bubbles is applied between the cathode and the anode.\nU.S. Pat. No. 7,067,204 to Nomura et al., describes an apparatus comprising an ultrasonic generator for creation of bubbles within a liquid, and a generator providing the excitation of electromagnetic waves in the liquid phase, for the implementation of the plasma discharge.\nJapanese Application JP2006273707 to Shibata et al. relates to the publication, “Synthesis of amorphous carbon nanoparticles and carbon-encapsulated metal nanoparticles in liquid benzene by an electric plasma discharge in ultrasonic cavitation field,” Ultrasonic Sonochemistry 13 (2006) 6-12, Institute of Multidisciplinary Research for Advanced Material (IMRAM), Tohoku University. This application illustrates a method and a device for producing a nanocarbon material that does not require an expensive production facility such as the ones normally required for dry treatment. It can easily produce the nanocarbon material because the application of high voltage is not needed and neither worsens nor deteriorates the working environment in a production premise, and at the same time considers safety factors. This method can remarkably reduce production costs by improving production efficiency because of its continuous production and recovery, and providing an alternative for mass productivity. The method comprises a process (A) for arranging electrodes, one cathode and one anode, connected to the power source; an ultrasonic horn connected to an ultrasonic generator within an organic solvent that fills a container; and a process (B) for generating an ultrasonic cavitation field by ultrasonic waves into the organic solvent, around the head of the ultrasonic horn; and effecting the thermal decomposition of the molecules in the organic solvent by applying a voltage to the electrodes so as to generate plasma discharge within the ultrasonic cavitation field adequate for the production of the nanocarbon material.\nU.S. Pat. No. 6,835,523 to Yamazaki et al. describes a “Method for fabricating with ultrasonic vibration a carbon coating,” which is a process for fabricating a carbon coating in a medium disposed on one side of an electrode connected to a high-frequency power supply. Ultrasonic vibrations are then supplied to the object.\nNone of the prior art, however, either individually or in combination, provides a method by which initiating and maintaining an abnormal glow volumetric sonoplasma discharge can be performed using a substantially lower amount of electrical power.\nThus there has existed a long-felt need for a method by which the sonoplasma discharge can be initiated and maintained with substantially less electrical power than is currently needed to accomplish the same result using the prior art. This is accomplished with this invention.\nThe current invention provides just such a solution by having a method and apparatus for initiating and maintaining an abnormal glow volumetric sonoplasma discharge (VSPD). With certain parameters of the electrical discharge and of the intensity of elastic vibrations, it is possible to initiate VSPD within a cavitating liquid medium. The mechanism for the initiation of VSPD is related to the breakdown of gas-phase microchannels formed by the growth cavitation bubbles. The method uses elastic vibrations (EV) in the frequency range 1,000-100,000 Hz with enough intensity for the development of cavitation phenomena; these vibrations are introduced into the liquid-phase working medium, and a source of direct, alternating (hertz and kilohertz range), high frequency (HF) (megahertz range) and ultrahigh frequency (UHF) (gigahertz range) electric field in liquid (DPS) provides the initiation and stable glow of VSPD. Resulting VSPD is characterized by volumetric glow in the frequency range of visible light and ultraviolet radiation in the entire cavitation-electric field, and is characterized by a rising volt-ampere characteristic curve.\nWhen a high-intensity ultrasonic field exceeding a cavitation threshold is induced within liquids, a new form of electric discharge is obtained, characterized by a volumetric glow electrical discharge throughout the space between the electrodes, having a rising volt-ampere characteristic curve that is inherent to abnormal glow discharge in gas. Such discharge within the liquid has the surface characteristic of micro bubbles, and can be used for the design of novel sonoplasma-chemical processes because of the extensive interface plasma. The heterogeneous liquid/gas-vapor system leads to a rise in diffusion rates of chemically active particles in the system and a more economical method to achieve the desired result(s)."}
{"text": "Over the last two decades, mass spectrometry has made tremendous strides in analyzing protein samples derived from a variety of different sample types. Coupled with electrospray ionization and various separation techniques, thousands of proteins may be identified and quantitated in a single sample. The most common approach used in the laboratory today involves some form of protein extraction followed by proteolytic digestion of protein sample of interest. The use of proteolytic enzymes like trypsin produces peptides that can easily be analyzed by a variety of different instrument configurations. This approach termed “bottom-up” proteomics, can be used to study the state of living cells as a function of their environment. One of the major advantages of the “bottom-up” approach is that the peptides produced have very similar physiochemical properties which makes for a straight forward separation of thousands of peptides in complex samples. Any separation approach coupled with tandem mass spectrometry can then be used to produce amino acid sequence information that is utilized to identify the proteins in a given sample. Although this technique is routine in many laboratories, there are limitations as to the amount of information that can be obtained when reducing intact proteins to their constituent peptides.\nIn contrast to “bottom-up” proteomics, “top-down” proteomics refers to methods of analysis in which protein samples are introduced intact into a mass spectrometer, without enzymatic, chemical or other means of digestion. Top-down analysis enables the study of the intact protein, allowing identification, primary structure determination and localization of post-translational modifications (PTMs) directly at the protein level. Top-down proteomic analysis typically consists of introducing an intact protein into the ionization source of a mass spectrometer, fragmenting the protein ions and measuring the mass-to-charge ratios and abundances of the various fragments so-generated. The resulting fragmentation is many times more complex than a peptide fragmentation, which may, in the absence of the methods taught herein, necessitate the use of a mass spectrometer with very high mass accuracy and resolution capability in order to interpret the fragmentation pattern with acceptable certainty. The interpretation generally includes comparing the observed fragmentation pattern to either a protein sequence database that includes compiled experimental fragmentation results generated from known samples or, alternatively, to theoretically predicted fragmentation patterns. For example, Liu et al. (“Top-Down Protein Identification/Characterization of a Priori Unknown Proteins via Ion Trap Collision-Induced Dissociation and Ion/Ion Reactions in a Quadrupole/Time-of-Flight Tandem Mass Spectrometer”, Anal. Chem. 2009, 81, 1433-1441) have described top-down protein identification and characterization of both modified and unmodified unknown proteins with masses up to ≈28 kDa\nAn advantage of a top-down analysis over a bottom-up analysis is that a protein may be identified directly, rather than inferred as is the case with peptides in a bottom-up analysis. Another advantage is that alternative forms of a protein, e.g. post-translational modifications and splice variants, may be identified. However, top-down analysis has a disadvantage when compared to a bottom-up analysis in that many proteins can be difficult to isolate and purify. Thus, each protein in an incompletely separated mixture can yield, upon mass spectrometric analysis, multiple ion species, each species corresponding to a different respective degree of protonation and a different respective charge state, and each such ion species can give rise to multiple isotopic variants.\nThe process of analyzing intact proteins in cell lysates by mass spectrometry (MS) is associated with a number of difficulties. Firstly, electrospray ionization (ESI) of protein mixtures from cell lysates can generate extremely complex mass spectra due to the presence of multiple proteins, each comprising its own charge state envelope, where each charge state envelope is the collection of mass spectral lines corresponding to plural charge states, and where each charge state correlates directly with the number of positively charged protons that are adducted to an otherwise charge-free molecule. Consequently, multiple charge state envelopes may be overlapping within any given mass-to-charge (m/z) range. In this example, multiple proteins overlap at the same m/z value that have different molecular weights and charges. Commonly used techniques in MS are often insufficient for simplifying these spectra because of the inherent peak overlapping as well as the inherent wide range of magnitudes of MS lines of ionized constituents, where such constituents may range from uninteresting small molecules to interfering biomolecules to the proteins of interest, themselves. Isolation of a specified charge state of a protein within such complex spectra does not typically alleviate the burden of multiple protein peaks overlapping, since the isolation of ions of a particular protein charge state will generally result in co-isolation of one or more additional ions. This co-isolation makes it a challenge not only to dissociate the protein in an attempt to identify it based on the fragments produced, but also to accurately determine the intact mass and sequence coverage of that protein.\nSo-called “front-end” separation techniques, such as liquid chromatography (LC) or ion mobility spectrometry (IMS), performed prior to introduction of samples into a mass spectrometer, may be implemented to reduce the overall complexity and provide an additional major benefit, which is the reduction of ionization competition at an ionization source. Unlike mixtures of proteolytic peptides typically analyzed in bottom-up experiments, intact proteins mixtures contain a wide range of molecular weights, isoelectric points, hydrophobicities, and other physiochemical properties that make it challenging to analyze these mixtures via any single separation technique in a comprehensive manner. Both of the above separation methods are associated with their own benefits and pitfalls. Liquid chromatography tends to require significant amounts of time per sample to separate individual proteins, although it is still common to have two or more proteins co-elute. Enhanced separation can reach the point of becoming more of “an art” than a standardized method, and the enhanced separation may be dependent on the user skill in the state-of-the-art. The latter technique, IMS, can rapidly separate certain proteins and/or charge states from others but IMS spectra are at least partially correlative with (i.e., not “orthogonal to”) mass spectra. The IMS method also suffers from ionization competition, requires extensive optimization and typically involves dynamic conditions to observe a full mass spectrum containing all charge states.\nProton transfer reactions, a type of ion-ion reaction that has been used extensively in biological applications for rapid separations of complex mixtures, addresses many of these aforementioned concerns. Experimentally, proton transfer is accomplished by causing multiply-positively-charged protein ions from a sample to react with introduced singly-charged reagent anions so as to reduce the charge of the multiply-charged protein ions. These reactions proceed with pseudo-first order reaction kinetics when the anions are present in large excess over the protein ion population. The rate of reaction is directly proportional to the square of charge of the protein ion (or other multiply-charged cation) multiplied by the charge on the anion. The same relationship holds for reactions of the opposite polarity as well. This produces a series of pseudo-first order consecutive reaction curves as defined by the starting multiply-charged protein ion population. Although the reactions are highly exothermic (in excess of 100 kcal/mol), proton transfer is an even-electron process performed in the presence of 1 mtorr of background gas (i.e. helium) and thus does not fragment the starting multiply-charged protein ion population. The collision gas serves to remove the excess energy on the microsecond time scale (108 collisions per second), thus preventing fragmentation of the resulting product ion population.\nProton transfer reactions (PTR) have been used successfully to identify individual proteins in mixtures of proteins. This mixture simplification process has been employed to determine charge state and molecular weights of high mass proteins. PTR has also been utilized for simplifying product ion spectra derived from the collisional-activation of multiply-charged precursor protein ions. Although PTR reduces the overall signal derived from multiply-charged protein ions, this is more than offset by the significant gain in signal-to-noise ratio of the resulting PTR product ions. The PTR process is 100% efficient leading to only single series of reaction products, and no side reaction products that require special interpretation and data analysis.\nVarious aspects of the application of PTR to the analysis of peptides, polypeptides and proteins have been described in the following documents: U.S. Pat. No. 7,749,769 B2 in the names of inventors Hunt et al., U.S. Patent Pre-Grant Publication No. 2012/0156707 A1 in the names of inventors Hartmer et al., U.S. Pre-Grant Publication No. 2012/0205531 A1 in the name of inventor Zabrouskov; McLuckey et al., Anal. Chem. 1998, 70:1198-1202; Stephenson et al., J. Am. Soc. Mass Spectrom. 1998, 8:637-644; Stephenson et al., J. Am. Chem. Soc. 1996, 118:7390-7397; McLuckey et al., Anal. Chem. 1995, 67:2493-2497; Stephenson et al., Anal. Chem. 1996, 68:4026-4032; Stephenson et al., J. Am. Soc. Mass Spectrom. 1998, 9:585-596; Stephenson et al., J. Mass Spectrom. 1998, 33:664-672; Stephenson et al., Anal. Chem., 1998, 70:3533-3544 and Scalf et al., Anal. Chem. 2000, 72:52-60. Various aspects of general ion/ion chemistry have been described in McLuckey and Stephenson, Mass Spec Reviews 1998, 17:369-407 and U.S. Pat. No. 7,550,718 B2 in the names of inventors McLuckey et al. Apparatus for performing PTR and for reducing ion charge states in mass spectrometers have been described in U.S. Pre-Grant Publication No. 2011/0114835 A1 in the names of inventors Chen et al., U.S. Pre-Grant Publication No. 2011/0189788 A1 in the names of inventors Brown et al., U.S. Pat. No. 8,283,626 B2 in the names of inventors Brown et al. and U.S. Pat. No. 7,518,108 B2 in the names of inventors Frey et al. Adaptation of PTR charge reduction techniques to detection and identification of organisms has been described by McLuckey et al. (“Electrospray/Ion Trap Mass Spectrometry for the Detection and Identification of Organisms”, Proc. First Joint Services Workshop on Biological Mass Spectrometry, Baltimore, Md., 28-30 Jul. 1997, 127-132).\nThe product ions produced by the PTR process can be accumulated into one or into several charge states by the use of a technique known as “ion parking”. Ion parking uses supplementary AC voltages to consolidate the PTR product ions formed from the original variously protonated ions of any given protein molecule into a particular charge state or states at particular m/z values during the reaction period. This technique can be used to concentrate the product ion signal into a single or limited number of charge states (and, consequently, into a single or a few respective mass-to-charge [m/z] values) for higher sensitivity detection or further manipulation using collisional-activation, ETD, or other ion manipulation techniques. Various aspects of ion parking have been described in U.S. Pat. No. 8,440,962 B2 in the name of inventor Le Blanc and in the following documents: McLuckey et al., Anal. Chem. 2002, 74:336-346; Reid et al., J. Am. Chem. Soc. 2002, 124:7353-7362; He et al., Anal. Chem. 2002, 74:4653-4661; Xia et al., J. Am. Soc. Mass. Spectrom. 2005, 16:71-81; Chrisman et al., Anal. Chem. 2005, 77:3411-3414 and Chrisman et al., Anal. Chem. 2006, 78:310-316.\nAnother difficulty associated with the mass spectrometric analysis of proteins in cell lysates by (MS) is that the fragmentation behavior for each charge state of a protein is generally unknown prior to the dissociation event. In particular, ions comprising some charge states can dissociate well while ions comprising other charge states may dissociate poorly. Isolation and dissociation of ions of a particular charge state therefore does not guarantee efficient dissociation or dissociation into a set of diagnostic fragments.\nA third challenge associated with intact protein analysis is the wide distribution of charge states produced for high molecular weight proteins typically in excess of 50 kDa. Here the starting signal can be divided into over 30 plus charge states, making tandem mass spectrometry of any given charge state produce a spectrum with low signal-to-noise ratio. The ability to produce ample sequence coverage for protein identification can therefore be difficult with a single tandem mass spectrum.\nA variety of ion activation (fragmentation) techniques can be used to produce structural information on intact proteins. The most commonly used approach termed collision-induced dissociation (CID) involves collisions of an isolated population of multiply-charged precursor ions with a neutral background gas. Most commonly, the multiply-charged precursor ions are accelerated using the fundamental frequency of motion of the defined ion population in order to collide with the neutral background gas so as to produce unimolecular dissociation events. This process leads to fragmentation along the amide backbone of the protein thus yielding amino acid sequence information. More extensive fragmentation of proteins can be obtained with higher collision energy processes termed HCD or high energy collision induced dissociation. Many times this involves multiple fragmentation events inside the collision cell thus producing more extensive sequence coverage. Another approach used to produce protein sequence coverage via ion activation is that of photodissociation (PD), where photons of a defined wavelength are used to excite the ion of interest. Two common types employed include ultra-violet (UV-PD) and infrared multiphoton dissociation (IRMPD). The latter is a high energy process where the rate of energy deposition in the ion far exceeds that of the dissociation process. Here fragmentation can be produced along any point in the protein backbone, or may yield amino acid side chain fragmentation as well. For IRMPD, this is a much lower energy process that is characterized by the presence of cleavages at amide bonds and losses of ammonia and water from the intact protein and fragment ions generated during irradiation. The time frame of the IRMPD experiment can be expanded to produce more extensive fragmentation as well. Ion-ion reactions using electron transfer reagent ions can also be employed as a fragmentation approach for intact proteins. Here an electron transfer event from the multiple-charged protein to the singly-charged anion produces backbone fragmentation of the protein with any posttranslational modifications still intact.\nTaken together, these ion activation approaches for tandem mass spectrometry produce many different complementary forms of fragmentation that can provide protein sequence information. Ideally, these approaches can be applied in a broad band fashion in order to increase sequence coverage of proteins and provide additional information on modifications, splice variants, and expression of single amino acid mutations. The application of these approaches in a broadband format (i.e. covering multiple charge states of the same intact protein) would provide a more comprehensive view of protein characterization and identification."}
{"text": "1. Technical Field\nThe present device is a system for protecting windows, doors, and similar openings to a building against hurricane-force winds, more particularly a system comprising specially designed wedge brackets, bolts or screws that project year-round from the corners the openings, and a protective panel of sufficient size to cover the opening.\n2. Background Information\nAn estimated six million single and double family dwellings, and many more commercial buildings, along United States coastlines are potentially affected by hurricanes as they spin their way up across the Atlantic Ocean or up through the Gulf of Mexico. As a hurricane bears down, homeowners and business people frequently prepare by nailing or screwing wooden boards up over windows and doors. Preparations are particularly feverish when a Category 4 or 5 hurricane is expected to make landfall. When a hurricane is expected to make landfall in the vicinity, the lines at lumber stores often stretch out into the parking lot and supplies of wooden boards, screws, and the like often run out. At the same time they are trying to protect their homes and valuables from the approaching hurricane, many people are simultaneously trying to pack up their children, elderly relatives, pets, clothing, photos, and other items to travel out of the path of the storm. They are concerned about getting out of the area before traffic jams occur. An easier, quicker system for protecting their homes and businesses is therefore much needed.\nThe present invention is a system for the homeowner or business person to use to protect their residences and other buildings from wind and rain damage. Anyone who is capable of lifting a board can implement this system over the windows and doors of a building."}
{"text": "There are two quite different types of electronic networks that are evolving: a standard telephone network and a data network (e.g., the Internet). The standard telephone network, such as a wireless telephony network and the POTS network, is designed to carry real-time messaging content. Capacity is allocated in real-time bandwidth and, once you have a connection of adequate bandwidth established between two points, data that is delayed is viewed as a network fault. Examples of communications that may be carried via a standard telephone network(s) include voice communications, multi-party conference calls, video conference calls, or the like.\nAnother characteristic of standard telephone networks is that each network is typically owned and controlled by a small number (typically a single company in the case of wireless networks) of large companies that historically have provided service directly to end users of the network and therefore have a billing type relationship with them. Where a connection is made through the facilities of another provider, there is typically a commercial contract in place between the two companies. These standard telephone networks, in part because of the relatively close relationship between the service vendors and network users, have the characteristic that the originator of a call can be readily identified, allowing “caller ID” service to be readily implemented and to be widely known and in fact expected.\nIn contrast, the second type of network (e.g., data networks such as the Internet, LANs, WANs, VPNs, and the like) were designed to move mostly one-way, non-real time data from point to point. In this type of network, the delay of data has typically not been regarded as a network fault. Additionally, some data networks, particularly the Internet, are far more disjoint than a standard telephone network. There are many more companies involved, and there is much less control of individual point-to-point end-user connections. It is typical that a company that provides transmission of data on the Internet has a tenuous commercial relationship with the originators of most of the data packets that it is carrying. In fact, Internet service providers (ISPs) protect the privacy and anonymity of their subscribers.\nThis tenuous commercial relationship with end users coupled with the relative ease with which the end-user computers that originate much of the traffic on the Internet can be anonymously enlisted in the service of third parties, leads to the fact that a “caller ID” type service is nearly impossible to implement on the Internet.\nIn recent years, these data networks have begun to evolve to provide real-time, two-way communications between parties. The communications may include, for example, voice-over-IP (VOIP), instant messaging, interactive video conferencing (e.g., web meetings), or the like.\nUsing the electronic data networks for real-time, two-way communications provide several advantages. In particular, using these electronic networks for real-time, two-way communications is relatively low cost and easily accessible. The proliferation of networks throughout today's society, particularly the Internet, has ensured ready access to a communications device capable of communicating with any other individual communicatively coupled to the same network. Essentially anyone with a computer, a personal data assistant, a wireless telephone, or the like can connect to the Internet and communicate with someone at a remote location within seconds. Likewise, companies can use internal networks (e.g., WANs, VPNs, or the like) to allow geographically dispersed employees to communicate in real-time using many of the same technologies. Notably, the communications can frequently occur with equipment already purchased as networks and access devices are generally already in place to handle data needs.\nAs this type of communications becomes more widespread, it will inevitably become a target for advertisers and telemarketers as a method to distribute advertising messages in vast quantities. Because of the low cost of distributing massive amounts of advertising, advertisers can economically transmit advertising communications with response rates that are orders of magnitude less than would be necessary to support more traditional means of advertising. Additionally, as discussed above, anonymity given the sender prevents “do not call” lists and caller-ID type mechanisms from providing an adequate solution.\nElectronic mail (e-mail) has already seen this problem. E-mail is a store-and-forward communications method in which one-way communications (as opposed to a two-way communications) are sent from one network node to another network node until the final destination is reached, where a recipient may or may not retrieve a message or respond. Because e-mail is inexpensive and advertisers can transmit massive amounts of e-mail quickly (and often automatically), e-mail advertisements (e.g., junk e-mail) are becoming a burden to networks and users alike. This use of e-mail to send massive amounts of advertisements is known as the e-mail “spam” problem.\nAttempts have been made to reduce the effect of e-mail spam on the end users. One such attempt is described in U.S. Pat. No. 6,052,709, wherein a system that attempts to filter incoming e-mail to identify junk e-mail is described. This system, however, only applies to e-mail, which, as described above, is a one-way communication, and does not apply to two-way, real-time communications, such as voice, video, real-time text, or the like.\nTherefore, there is a need for a method and system to identify and filter unsolicited real-time, two-way communications."}
{"text": "The semiconductor integrated circuit (IC) industry has experienced rapid growth. Technological advances in IC materials and design have produced generations of ICs where each generation has smaller and more complex circuits than the previous generation. However, these advances have increased the complexity of processing and manufacturing ICs and, for these advances to be realized, similar developments in IC processing and manufacturing are needed.\nIn the course of IC evolution, functional density (i.e., the number of interconnected devices per chip area) has generally increased while geometry size (i.e., the smallest component (or line) that can be created using a fabrication process) has decreased. This scaling down process generally provides benefits by increasing production efficiency and lowering associated costs. Such scaling-down has also significantly decreased the space separating contact openings (or windows) from adjacent devices of ICs. Contact openings provide contact between various devices and features of the integrated circuit. Due to the scaled down devices and decreased space separation between devices, it has been observed that conventional processing provides a smaller than desirable contact process window, which leads to restrictive processing and design issues. For example, the smaller contact process window results in design rules requiring a minimum spacing between the contact openings and device features (e.g., gate structures), which provides a smaller than desirable margin of contact/gate structure overlay. Further, if the minimum spacing between the contact openings and such device features varies, poor device performance results, such as contact/gate structure short and contact open issues.\nAccordingly, what is needed is a method for making a semiconductor device that addresses the above stated issues."}
{"text": "In hierarchical computer storage systems, fast and intensively used storage are paired with arrays of slower and less frequently accessed data devices. One example of high-speed, expensive memory is a direct access storage device file buffer (DASD). Slower storage devices include tape drives and disk drive arrays, which are less expensive than a DASD.\nOne such hierarchical storage system is a virtual tape storage system. Such a virtual tape storage system may include, for example, one or more virtual tape servers (“VTS”) in combination with one or more data storage and retrieval systems, such as the IBM TotalStorage® 3494 Enterprise Tape Library. During operation, each virtual tape storage system is communicating data from one or more hosts, and is providing data to a second VTS for copying.\nData disaster recovery solutions include various “peer-to-peer” copy routines where data is backed-up not only remotely, but also continuously (either synchronously or asynchronously). In order to communicate duplexed data from one host processor to another host processor, or from one storage controller to another storage controller, or some combination thereof, a substantial amount of control data is required for realizing the process. A high overhead, however, can interfere with a secondary site's ability to keep up with a primary site's processing, thus threatening the ability of the secondary site to be able to recover the primary in the event a disaster occurs.\nDisaster recovery protection for the typical data processing system requires that primary data stored on primary DASDs be backed-up at a secondary or remote location. The physical distance separating the primary and secondary locations can be set depending upon the level of risk acceptable to the user, and can vary from several kilometers to thousands of kilometers.\nUsing prior art methods, in the case where, if the peer-to-peer subsystems, i.e. both virtual tape servers, are shutdown for normal service, and for some reason only one of those virtual tape servers becomes operational, then the peer-to-peer cluster must wait until both tape servers are again operational before going online to the host computer. Therefore using these prior art methods, if a second virtual tape server fails while the first virtual tape server is shutdown for maintenance, then the entire peer-to-peer system becomes unavailable until both virtual tape servers are again operational.\nWhat is needed is a method to distribute information about the status of a peer-to-peer data storage system across a plurality of system components such that the system itself can use that stored system information to return to operation even if all the virtual tape servers are not operational."}
{"text": "This disclosure is related to the field of neutron well logging measurements for determining petrophysical properties of subsurface formations traversed by a wellbore. More specifically, the disclosure relates to methods and apparatus for determining whether signals generated by a radiation detector in a well logging instrument were induced by neutrons or gamma rays. The disclosure also relates to methods and apparatus for gain stabilization of radiation detectors.\nWell logging instruments known in the art include various types of radiation detectors that are capable of detecting neutrons and gamma rays. Radiation detectors include gas-filled tubes in which the gas becomes ionized following a radiation event in the tube that the detector is configured to detect. Other types of radiation detectors include scintillation detectors, which may comprise a radiation sensitive scintillation crystal optically coupled to a photomultiplier tube. Scintillation detectors may be used, for example and without limitation, detecting gamma rays so as to be able to characterize the energy of the detected gamma rays by measuring the amplitude of each signal pulse generated by the photomultiplier tube.\nThere remains a need for a radiation detector, or a method or system utilizing a radiation detector, that delivers improved well logging performance."}
{"text": "Cellular communications systems are well known in the art. In a typical cellular communications system, a geographic area is divided into a series of regions that are referred to as “cells,” and each cell is served by one or more base stations. A base station may include baseband equipment, radios and antennas that are configured to provide two-way radio frequency (“RF”) communications with mobile subscribers that are geographically positioned within the cell. A common cellular communications system network plan involves a base station serving a cell using three base station antennas, wherein each base station antenna serves a 120 degree “sector” of the cell in the azimuth plane. The base station antennas are often mounted on a tower or other raised structure, with the radiation pattern (“antenna beam”) that is generated by each base station antenna directed outwardly to serve the respective sector. Typically, a base station antenna is implemented as a phase-controlled array of radiating elements, with the radiating elements arranged in one or more vertical columns. Herein, “vertical” refers to a direction that is perpendicular relative to the plane defined by the horizon.\nAs demand has grown for cellular communications systems to support increased capacity and provide enhanced capabilities, a variety of new cellular services have been introduced. These new services typically operate in different frequency bands from existing services to avoid interference. When new services are introduced, the existing “legacy” services typically must be maintained to support legacy mobile devices. Thus, as new services are introduced, either new cellular base stations must be deployed or existing cellular base stations must be upgraded to support the new services in the new frequency bands. In order to reduce cost and the total number of base station antennas deployed, base station antennas are now available that include at least two different arrays of radiating elements, where each array of radiating elements supports a different type of cellular service in a different frequency band. Such antennas are typically referred to as multi-band antennas."}
{"text": "1. Field of the Invention\nThis invention relates to the recovery of proteinaceous matter from such waste protein-containing liquids as whey and tannery unhairing waste.\n2. Description of the Prior Art\nPrecipitation of proteins from whey by heat denaturation is an old and straightforward procedure. In fact, the various whey proteins have been shown to be heat-denatured at different rates (J. Dairy Sci. 38, 351, 1955). However, the major whey proteins, .alpha.-lactalbumin and .beta.-lactoglobulin are rather resistant to denaturation by heat at the pH and concentration at which they occur in acid whey (J. Dairy Sci. 49, 694, 1966; J. Dairy Res. 37, 233, 1970). Both of these proteins are stabilized in their native conformations by internal disulfide bonds (Biochem. Biophys. Acta 200, 184, 1970; J. Dairy Sci. 57, 1152, 1974) and, at least in the case of .beta.-lactoglobulin, irreversible denaturation and the simultaneous appearance of new sulfhydryl groups has been thoroughly established (J. Dairy Sci. 33, 890, 1950; JACS 79, 126, 1957; J. Dairy Sci. 52, 585, 1969).\nTannery unhairing waste is an opaque, noxious liquid that has a 5-day biological oxygen demand (BOD) of over 19,000 ppm (JALCA 69, 50, 1974). Methods of disposing of this waste, which contains quantities of proteinaceous materials, are the subject of world wide studies (J. Water Pollut. Contr. Fed. 44, 1080, 1972; ibid 43, 998, 1971)."}
{"text": "1. Field of the Invention\nThis invention relates to punch presses, and more specifically to gaging used to locate a workpiece with respect to the die and hence with respect to a previously prepared program of punch press operation.\n2. Prior Art\nVarious types of gaging have been provided on punch presses to orient a workpiece with respect to a program of press operations. The problem is to get an exact precise dimension of predetermined size between the gaging surface and the centerline of the die. For example, at typical dimension is 48.000 inches plus or minus 0.0005 inch. Thus with such a dimension, the entire range of permissible locating errors is only one part in 48,000. During usage, wear becomes an immediate problem necessitating removal, grinding and shimming. Further, prior X-axis gaging has been so constructed as to be virtually specialized for a specific punch press construction."}
{"text": "This invention relates to electrical connectors and more particularly to a female blade terminal comprising a part of an electrical connector.\nElectrical connectors comprising interengaging male and female terminals are well known in the art and are in wide use for establishing electrical connection between various electrical components. One popular and widely used electrical connector comprises a male blade terminal adapted to be slidably received in a cage or socket defined by a female blade terminal.\nThere are many design parameters that must be considered when designing a connector of the blade type. Specifically, the connector must have a simple and inexpensive design; must have a relatively small size; must provide maximum current carrying capacity; must provide a low insertion force so as to allow easy connection and disconnection of the terminal elements; must provide sufficient grasping action as between the terminal elements to preclude inadvertent separation of the elements; must provide reliable and consistent operation for sustained periods of time; and must discourage corrosion and contamination of the electrical interface provided by the connector.\nWhereas a myriad of electrical connectors of the blade type have been proposed and are in common usage, none of the prior art blade type connectors totally satisfy all of the noted design parameters."}
{"text": "1. Field of the Invention\nThe present invention relates to a receiving apparatus that receives radio image signals transmitted from a body-insertable apparatus, such as a swallowable capsule endoscope, inserted into a subject body by employing plural antennas arranged outside of the subject body. More particularly, the present invention relates to the receiving apparatus that performs image processing on the received radio image signals.\n2. Description of the Related Art\nIn recent years, a capsule endoscope having an imaging function and a radio communication function makes an appearance in a field of endoscope. The capsule endoscope is swallowed from a mouth of a patient, i.e., the subject body for an observation (examination), and is eventually discharged naturally from a living body (human body) of the patient. While the capsule endoscope is inside the subject body, i.e., during an observation period, the capsule endoscope travels through organs (through body cavity) such as a stomach and a small intestine while following peristaltic motion of the organs, and the capsule endoscope sequentially images inside the organs by using the imaging function.\nFurther, during the observation period, i.e., while the capsule endoscope travels through the organs, image data obtained inside the body cavity by the capsule endoscope are sequentially transmitted outside of the subject body by the radio communication function such as Bluetooth, and stored in a memory that is provided inside an external receiving apparatus. When the patient carries around the receiving apparatus that has the radio communication function and a memory function, the patient can freely move without any inconveniences even during the observation period, i.e., after swallowing the capsule endoscope until discharging the capsule endoscope. After the observation, a doctor or a nurse can display an image inside the body cavity on a display unit such as a display based on the image data stored in the memory of the receiving apparatus to make a diagnosis.\nIn general, the receiving apparatus includes plural antennas dispersively arranged outside of the subject body to receive image signals transmitted from the capsule endoscope, and switchably selects one antenna that has small error in receiving the image signals to receive the image signals, and performs image processing on the received image signals. In a medical apparatus proposed in Japanese Patent Application Laid-Open No. 2003-325439, a capsule ID corresponding to a unique number such as a capsule serial number is superposed on image signals of the capsule, and externally transmitted in a frame format, so that a capsule-side switch for identifier setting is not necessary and an amount of information of the identifier is reduced."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for the assembly of and the interconnection by diffusion of bodies of metal alloys. More particularly the invention relates to methods of the kind comprising a sequence of activation of parts to be joined and a sequence of applying pressure when hot which gives rise to interconnection by interdiffusion of the solids.\n2. Summary of the Prior Art\nMetal alloys of which interconnection can be effected by these methods are alloys of the kind which are resistant to heat or to corrosion comprising, having in proportions by weight, at least 50% of a metal of the group nickel, cobalt and iron or an alloy of at least two of these metals. Certain of these alloys are designated by the term \"super-alloys\" and in such alloys the remainder is constituted by elements such as chromium, aluminum, molybdenum, titanium, tungsten, niobium, etc., which form a solid solution and form intermetallic compounds or alternatively form with carbon dispersed phases providing the required functional characteristics (mechanical strength and resistance to corrosion) at elevated temperatures. These alloys are used, for example, in the construction of turbo machines.\nBy \"body\", there is to be understood herein and in that which follows, simultaneously, and indifferently, either relatively large parts of predetermined shape and dimensions, for example intended to form a welded assembly, or powder grains intended for the production of parts by sintering and by compaction when hot. In practice, the interconnections, junctions or joints envisaged by the invention concern simultaneously, either two relatively large parts, or such a part and a powder, or powders, and, for the interconnection or joining of two relatively large parts, the invention is applicable whether or not an external deposit element is incorporated and which is constituted, for example, by a brazing compound. The assembly method and interconnection concerned thus encompasses different assembly and connection methods either by welding/diffusion, or by brazing/diffusion or by sintering or by compaction when hot. Briefly, the principles of these methods generally known to the man skilled in the art are will now be summarized.\nThe method of welding/diffusion involves welding the solid phase in which the parts are kept in contact under a given pressure and brought to a predetermined temperature, lower than the initiation of the melting temperature, over a predetermined time period. These operational conditions lead to local plastic deformations of the contact surfaces, which, in turn, give rise to an intimate contact of the latter thus rendering possible the migration of atoms between elements and/or recyrstallisations at the interface. In the ideal case, once the operation is completed no presently available technique enables the initial contact interface to be distinguished, whether this is considered from the microphotographic, chemical or mechanical aspects.\nAs for brazing/diffusion, it consists in the insertion between the two parts to be assembled of a thin layer (foil or powder) of an alloy particularly with a nickel base comprising additions (fluxes) such that the temperature of the liquids of the alloy will be less than the temperature of the initiation of melting of the superalloy. Heating causes initially the fusion of this layer and its connection with the superficial contact layers of the parts, then the reduction in the local proportion of the or each flux by migration into adjacent zones of the parts. In fact, it does not relate to a welding method by solids interdiffusion but the application of a known brazing technique then the diffusion of the fluxes of the brazing material in order to elevate the temperature of fusion of the connecting layer.\nSintering is likewise shown to be of interest for the production of parts of super-alloys, given the difficulties encountered during the use of conventional methods such as machining or forging. Sintering enables the production of complex shaped parts starting from a superalloy powder which, after having been shaped, is submitted to the operation of sintering itself. The latter consists either in heating the powders to a predetermined temperature after compaction when cold in a mould or matrix, or to place the powders under pressure at a predetermined temperature, lower than the temperature of commencement of melting fusion or the latter enables the production of a consolidation of the part by a phenomenon analagous to welding by diffusion at the zone of each grain. The physical and mechanical properties of the sintered parts clearly depend upon the quality of the junctions or connections formed between the powdered grains.\nThe assembly of the bodies can be carried out either between similar materials, or between different materials by their composition or their nature: a dense material and a powdered material for example.\nThe superior mechanical characteristics of the alloy used, are accompanied on the negative side by difficulties of fabrication of the parts or of the assemblies, especially by sintering or diffusion welding.\nIt is in practice apparent during tests of the mechanical strength of assemblies of parts of alloys welded by diffusion, and micrographic examination, that for certain grades of super-alloys or for certain special alloys of the non-oxidizing type, the connection produced after welding by diffusion is of mediocre quality; similarly the strength of sintered parts produced by super-alloys is poor, the cohesion of the assembly being less than values conventionally obtained. In certain cases no joint or connection is possible.\nObservations based on electron microscope pictures, and the study of constituents by microanalysis have enabled evidence to be produced that these defects were due to the formation in the superficial contact layers of segregations or precipitates which may harm and even prevent diffusion. These segregations or precipitations result from the migration of certain constituents during the increase in temperature preceding the operation which effects the joint or connection.\nThis phenomenon of the formation of a barrier to diffusion is more particularly apparent for super-alloys containing, in addition to a non-negligeable amount of carbon, a relatively high quantity of titanium, for example upwards of 0.15% of carbon and 1.5% of titanium. It has been shown in this case, after diffusion welding, for example on one hand the presence in the interface of a quasi-continuous boundary formed by segregation of titanium compounds comprising especially carbon and titanium and, on the other hand, the absence of recrystallisation. FIG. 1a shows a microphotograph, obtained at an enlargement of 2500 times, at one example of an assembly of parts in which the phenomena hereinbefore noted may be observed. Mechanical tests confirm that the joint or connection thus produced is defective.\nIn French Patent Specification No. 2 380 354, with the object of resolving these difficulties, a preparatory treatment before the connection operation is proposed, consisting in a roughening treatment by heating in an enclosure where a hydro-halogen atmosphere is circulated composed of a mixture of hydrogen and of a hydrogen halide. French Patent Specification Nos. 1 170 557 and 1 243 415 also set out the conditions for carrying out the treatment in a fluoride atmosphere produced from chromium fluoride and ammonium fluoride in the form of a cement in the treatment enclosure obtained under a reducing atmosphere, for example of hydrogen. The use of the atmosphere of hydrogen can, however, lead to certain difficulties and constraints particularly the way of putting the technique into practice."}
{"text": "Prior art cardiac monitoring devices such as cardioverter defibrillators, holter monitors, or ICU monitors often use the same criteria to detect and subclassify cardiac arrhythmias. Prior art devices commonly utilize a single criterion such as cycle length to detect tachyarrhythmias. These devices generally measure cardiac cycle length by measuring the time between the large electrical deflections produced when the ventricles depolarize. The electrical deflections are sensed when the signal amplitude exceeds the amplitude of a programmed threshold. Prior art devices detect tachyarrhythmias by determining when the cardiac cycle length, or time between consecutive ventricular contractions, falls below a programmed level. The programmed levels are typically as follows: a cycle length greater than 500 milliseconds (ms) is identified as normal, a cycle length between 500 and 333 ms is classified as monomorphic tachycardia, and a cycle length less than 333 ms is identified as polymorphic tachycardia.\nThus the effectiveness of prior art devices of this type depend on the accuracy of cycle length detection. There are two major disadvantages with these devices. First, these detectors sometimes miss low amplitude electrical activity, such as during ventricular fibrillation, which can cause the detector to miss dangerous polymorphic ventricular tachycardias. Second, cycle length is not an effective discriminator between monomorphic and polymorphic arrhythmias even when the threshold detectors sense the electrical events appropriately.\nOther prior art devices have been developed to overcome these disadvantages by utilizing additional parameters which can be programmed by a physician. Examples of some detection parameters include cycle length cutoff for monomorphic ventricular tachycardias, cycle length cutoff for polymorphic ventricular tachycardias, cycle length regularity, and QRS width. The disadvantage with these programmable devices is that a large number of detection parameters must be programmed by a physician. Further, this programming process can be complex, time consuming and prone to physician error.\nIt is highly desirable to have a non-programmable device and method for detecting tachyarrhythmias, i.e., a device which does not require programming by a physician. Further, a device is needed for use with cardioverter defibrillators, or monitors which can more accurately classify and discriminate between normal, monomorphic and polymorphic arrhythmias. Still further yet, it is desirable to have a device which can accurately sub-classify and discriminate between different types of arrhythmias within the monomorphic arrhythmia class or the polymorphic arrhythmia class. Thus, it is highly desirable to have an improved device and method which eliminates the prior art problems of missing dangerous polymorphic ventricular tachycardias when there is low electrical activity, and discriminating between monomorphic and polymorphic arrhythmias."}
{"text": "Conventionally, as disclosed in a patent document, Japanese Patent Laid-Open No. JP H10-19156 A (Patent document 1) listed below, a control unit performs a feedback control for controlling an electric current flowing in a solenoid of an electromagnetic valve.\nThe control unit described above includes: a switching section (i.e., a Pulse Width Modulation (PWM) circuit) on a power supply path toward the solenoid for flowing an electric current to the solenoid when being turned ON; a detector (i.e., an electric current detection circuit) for detecting an actual electric current value flowing in the solenoid; and a feedback controller (i.e., a microcomputer) for setting a duty ratio, which allows the detected actual electric current value to follow a target electric current value, and generating a PWM signal having the set duty ratio in a preset cycle for supplying the signal to the switching section.\nThe oil pressure valve, which is operated by the solenoid, is provided in the hydraulic circuit of the automatic transmission. Therefore, the control unit disclosed in the patent document 1 is used for the control of the automatic transmission. In recent years, the hydraulic circuit has a complicated structure, and a configuration of such structure includes two or more hydraulic valves in the oil circulation portion of the circuit, among which one or more hydraulic valves may be operated by the solenoid.\nIn such a configuration, a coupled oscillation may occur, which results from an oil pressure effect bouncing around between the multiple hydraulic valves. The coupled oscillation in such a configuration/structure is confirmed by the inventor of the present application. The mechanism of how coupled oscillation occurs in the hydraulic circuit is understood as follows.\nThe propagation rate of oil pressure affects the characteristic of the coupled oscillation such as frequency, amplitude and the like. The propagation rate is determined by the viscosity of the oil, and the viscosity of the oil changes according to the oxidization of the oil and the environmental temperature of the oil in which the oil is used. Therefore, when the viscosity of the oil changes according to the change of the environmental temperature of the oil, for example, the oscillation of the oil in the circuit may become noticeable (i.e., the oscillation exceeding an allowable level has occurred), thereby coupling the oscillation of many parts of the oil and the circuit to result in the coupled oscillation.\nFor example, when the hydraulic circuit has three hydraulic valves in the circulation portion, an influence of the operation of the first oil pressure valve is transmitted to the second oil pressure valve through the oil. Therefore, an input pressure of the second oil pressure valve is not stabilized, and a valve position of the second oil pressure valve is not converged (i.e., is not stabilized). Further, an influence of the operation of the second oil pressure valve is transmitted to the third oil pressure valve through the oil. Therefore, an input pressure of the third oil pressure valve is not stabilized, and a valve position of the third oil pressure valve is also not converged. Furthermore, an influence of the operation of the third oil pressure valve is transmitted to the first oil pressure valve through the oil. Therefore, an input pressure of the first oil pressure valve is not stabilized, and a valve position of the first oil pressure valve is not converged.\nThus, the coupled oscillation occurs from the coupling of the effects from each of the hydraulic valves, which is understood as causing a continuous operation of the same valve. That is, as the convergence of the valve position of each of the hydraulic valves stays unachieved for a long time (i.e., the continuous operation of the valve lingers on), and the oil pressure does not really attenuate, causing a continuation of the oscillation of the oil. In such a situation, the controllability of the automatic transmission may deteriorate."}
{"text": "DE application 10222192 describes pentafluorosulfanylbenzoylguanidines as NHE1 inhibitors. The processes described therein for preparing these compounds, however, result in low yield and require reagents and reaction-conditions that necessitate great technical complexity or are unsuitable for preparation on a relatively large scale. It has now been found that said disadvantages can be avoided by a novel efficient synthesis which starts from commercially available 4-nitrophenylsulfur pentafluoride."}
{"text": "A programmable device (e.g., a programmable microcontroller) contains configuration registers which hold configuration data to establish functional blocks (e.g., which perform user-defined logic functions), I/O blocks (e.g., which configures input/output blocks interfacing to external devices), and/or signal routing resources (e.g., which connect the functional blocks to each other and/or the I/O blocks).\nThe configuration data may be represented as configuration bits of configuration registers (e.g., stored as volatile memory). Upon the boot-up of the programmable device, the configuration data stored in non-volatile memory may be copied to the configuration registers of the volatile memory. However, the configuration data residing in the configuration registers may be compromised due to several factors.\nFor example, an unintended software execution may create a write-over condition where improper data may be written to the configuration registers. Additionally, cosmic rays, X-rays, and/or other environmental factors may cause the configuration data to degrade (e.g., flip); these are known as soft errors. These errors (e.g., the write-over, soft errors, etc.) of the configuration data may compromise the functional block, the I/O blocks, and/or the routing resources, thereby rendering the programmable device inoperable for its intended purposes.\nIn case when the programmable device is used in a critical condition (e.g., involving an emergency situation) or life critical function, the reliability of the programmable device becomes ever more critical. For instance, the functionality of an airbag deployment system may rely on the operation of a programmable device. Furthermore, the configuration data (e.g., bits) become more susceptible to the soft errors as the feature size (e.g., a silicon geometry) of the programmable device gets smaller."}
{"text": "1. Field of the Invention\nThe present invention relates to a data recording apparatus for recording data on a record (recording, recordable or recorded) medium, a method therefor, a data reproducing apparatus for reproducing data recorded on a record medium, a method therefor and a record medium on which data has been recorded. More particularly, the present invention relates to an information providing/collecting apparatus for providing and collecting so-called multimedia information, such as video information and music information, or program information and a method therefor.\n2. Related Background Art\nAs a data record medium on which information signals, such as Audio data, video data and various data items, are recorded, means for optically recording information signals, specifically, a so-called compact disk (CD) for use in the music field and a CD-ROM which meets the CD standard and which is used for data have been used all over the world in recent years.\nHitherto, information providing service has been realized as a so-called data base system and a personal computer communication system in each of which a user terminal (a terminal of an information collecting side) and an information provider are connected to each other through, for example, the telephone line to enable information required by the user to be taken out. Another information providing service has been realized with which a large-capacity medium, such as a so-called CD-ROM having encoded information recorded thereon is distributed and key information for decoding encoded information is transmitted to the user by, for example, communication so that encoded information recorded on the CD-ROM is decoded and decoded information is copied on a hard disk or the like so as to be used.\nMoreover, a technique has been disclosed in Japanese Patent Publication No. 2-60007 in which a password formed by encoding a file key by using a code key is supplied to a computer; and a program written on the record medium is decoded by a coding mechanism to prevent copying and sharing of the software program.\nHitherto, all of information items recorded on the foregoing CD or the CD-ROM are read by a reproducing apparatus and copied onto, for example, a hard disk. Then, data copied onto the hard disk is supplied to an encoder system for the CD or the CD-ROM to newly make a CD or a CD-ROM so that a pirate edition is easily manufactured. As described above, the security function, such as the copy protection, has been unsatisfactory.\nThe foregoing problem is also critical for a so-called digital video disk (DVD), which is expected to be a data record medium for a next generation.\nOn the other hand, in the conventional information providing service, a method has been employed in which key information for decoding is transmitted to a user in such a manner that key information is transmitted by means of voice through a telephone line. Thus, key information has not been encoded particularly. However, the foregoing method has a risk in view of keeping security.\nIn the case where communication is employed to transmit key information, one-to-one connection is usually established. Therefore, there is substantially no risk of key information being stolen. However, in the case where key information is transmitted through a network, there arises a problem in protecting key information.\nTherefore, in an information providing system, in which mediums, on each of which encoded information has been recorded in a large quantity, are distributed by the information provider; and only in a case where a user requires information to obtain from the medium, key information for decoding the code is supplied and accounting is performed, the problem in view of security when key information is transmitted results in a risk to arise in that key information can be obtained by a person except the subject user. In the foregoing case, the information providing system cannot be held. If whether or not the user is a formal user cannot be specified, there is a risk that account is put down to'another person. Also in the foregoing case, the information providing system cannot be held.\nThus, security improvement in transmitting key information from an information provider to a user and reliable specification of a user are important requirements."}
{"text": "1. Field of the Invention\nThe invention relates to the protection of integrated circuits from electrostatic discharge (ESD), and more particularly to the protection of NMOS transistors by an embedded parasitic silicon controlled rectifier (SCR) which triggers at a very low voltage.\n2. Description of the Related Art\nThe protection of integrated circuits from electrostatic discharge (ESD) is a subject which has received a lot of attention from circuit designers because of the serious damage that ESD can wreak as device dimensions are reduced. Workers in the field and inventors have proposed many solutions, many trying to solve the problem of protecting sub-micron (1 micron=10xe2x88x926 meter) devices while still allowing them to function unencumbered and without undue, or zero, increase of silicon real estate. The main thrust of ESD protection for MOS devices is focused on the use of parasitic npn and pnp bipolar transistors, which together form a silicon controlled rectifier (SCR). Unwanted as this SCR normally is, it can safely discharge dangerous ESD voltages as long as its trigger voltage is low enough to prevent gate oxide breakdown of the MOS devices of which it is a part.\nAmong ESD protection devices SCR protection shows good clamping capability (a very low holding voltage) compared to Gated-NMOS, however, a larger trigger voltage and latch-up concerns are always the drawbacks. See below a discussion of the graph of FIG. 1. In the sub-micron technologies, the Gated-NMOS is mostly used for a robust ESD design. A problem exists for the 0.15 micron process because the Gated-NMOS snapback voltage (larger than 5 volt) is higher than the gate oxide breakdown of 10 Million volt/cm. For 0.6 and 0.5 micron high voltage technology, the Gated-NMOS will generally show its weakness on high voltage (12 volt, 40 volt) ESD due to higher drain junction breakdown and higher snapback voltage.\nThe following publications describe low voltage lateral SCR structures to protect the input and output circuitry of an integrated circuit during an ESD event:\nxe2x80x9cA Low-Voltage Triggering SCR for On-Chip ESD Protection at Output and Input Pads,xe2x80x9d A. Chatterjee and T. Polgreen, IEEE Electron Device Letters, Vol. 12, No. 1, Jan., 1991.\nxe2x80x9cLateral SCR Devices with Low-Voltage High-Current Triggering Characteristics for Output ESD Protection in Submicron CMOS Technology,xe2x80x9d Ker, IEEE Transactions On Electron Devices, Vol.45, No.4, April 1998, pp.849-860.\nFIG. 1 is a graph of the I-V characteristics of a 0.15 micron process SCR/Gated-NMOS device. Curve 1 (dotted line) shows the characteristics for the SCR and Curve 2 (solid line) shows the characteristics for the Gated-NMOS. It is obvious that the SCR has the lower holding voltage while the Gated-NMOS has the lower trigger voltage. Another drawback of the SCR\"\"s is the latch-up concern.\nWe now describe in FIG. 2 a prior art low-trigger SCR for on chip ESD (IEEE, Electron Device Letters, Vol. 12, No. 1, January, 1991, see reference above). In a substrate 21 is embedded an n-well 32. Also implanted in substrate 21 is gated NMOS T1, comprising n+ drain 22, gate 23, and n+ source 24. Implanted next to n+ source 24 is p+ diffusion 25. Implants 23, 24 and 25 are tied to a voltage reference 39 (typically ground). N+ drain 22 is connected to chip pad 38. Implanted in the n-well are n+ diffusion 26 and p+ diffusion 27 which also are connected to chip pad 38. Halfway straddling the n-well is n+ diffusion 28 which is the drain of gated NMOS T2. T2 further comprises gate 29 and n+ diffusion (source) 30. A p+ diffusion 31 is implanted next to n+ diffusion 30. 29, 30, and 31 are connected to 39. Therefore, both T1 and T2 are grounded gate NMOS transistors. 26, 27, 28, 30, 31, and 32 make up the external SCR (external because the SCR is largely external to the gated NMOS transistors T1, T3, and T4). The latter two transistors are not further described since they follow the pattern of T1. The SCR itself comprises a parasitic bipolar pnp transistor Q1 and a parasitic bipolar npn transistor Q2. Drawbacks of this layout are a) the external SCR and b) low current capacity resulting in poor protection efficiency.\nFIG. 3 is a schematic diagram of the layout of FIG. 2 where the same numbers indicates the same items.\nAnother prior art lateral low-voltage, high-current SCR (IEEE Transactions On Electron Devices, Vol.45, No.4, April 1998, see reference above) is shown in FIG. 4 which is described in the cited IEEE document on page 851. Drawbacks of this design are:\na) poor I/O and VDD/GND connections causing latch-up concerns,\nb) the SCR is still outside of the main protection area leading to poor efficiency. Latchup occurred, see Arrow A, at a voltage difference of larger than 0.7 volt caused by diode Dp2 because of P-N diode forward cut-in.\nOther related art is described in the following U.S. Patents:\nU.S. Pat. No. 5,872,379 (Lee) describes a low voltage turn-on SCR to provide protection to the input and output circuitry of an integrated circuit during an ESD event.\nU.S. Pat. No. 5,907,462 (Chatterjee et. al) teaches a gate coupled SCR, where the stress voltage is coupled from a pad to a gate electrode causing a NMOS transistor to conduct, thus triggering the SCR.\nU.S. Pat. No. 5,939,756 (Lee) discloses an added P-well implantation for uniform current distribution during an ESD event to provide improved protection to the input and output circuitry of an integrated circuit.\nIt should be noted that none of the above-cited examples of the related art show part of the drain and the p+ diffusion in the n-well, nor having the drain connection tightly tied together at the p+ diffusion and the n+ drain as in the presently disclosed invention.\nIt is an object of the present invention to provide circuits and methods for an embedded SCR implemented in either 5 volt or 12 volt I/O devices.\nAnother object of the present invention is to provide for an embedded SCR which operates at less than or equal to 2 volts to prevent gate oxide damage, particularly in 0.15 micron 5 volt technology.\nA further object of the present invention is to provide for an embedded SCR where the ESD pass voltage of 8,000 volt is achieved for the 12 volt process.\nYet another object of the present invention is to provide a latch-up free SCR.\nThese objects have been achieved by inserting the p+ diffusion and the n-well in the drain side and a part of the drain to form a low-trigger, high efficiency SCR. Further, the layout is such that the drain connection is tightly tied together at the p+ diffusion and the n+ drain making that connection very short and thereby preventing latch-up. The parasitic SCR is made more efficient because it is entirely within the n+ diffusions (the source of the grounded gate NMOS transistor) at either side of the structure and, thus, called an embedded SCR. For the 12 volt I/O device the design is modified by placing each of two n+ drains in its own n-type doped drain (ndd) area and straddling halfway the n-well. The structure is repeated as required and a p+ diffusion is implanted at both perimeters and connected to the nearest n+ source and a reference voltage."}
{"text": "The present invention relates to ether compounds and pharmaceutically acceptable salts thereof; methods for synthesizing the ether compounds; compositions comprising an ether compound or a pharmaceutically acceptable salt thereof; and methods for treating or preventing a disease or disorder selected from the group consisting of a cardiovascular disease, dyslipidemia, dyslipoproteinemia, a disorder of glucose metabolism, Alzheimer\"\"s Disease, Syndrome X, a peroxisome proliferator activated receptor-associated disorder, septicemia, a thrombotic disorder, obesity, pancreatitis, hypertension, renal disease, cancer, inflammation, and impotence, comprising administering a therapeutically effective amount of a composition comprising an ether compound or a pharmaceutically acceptable salt thereof. The ether compounds and compositions of the invention may also be used to reduce the fat content of meat in livestock and reduce the cholesterol content of eggs.\nObesity, hyperlipidemia, and diabetes have been shown to play a casual role in atherosclerotic cardiovascular diseases, which currently account for a considerable proportion of morbidity in Western society. Further, one human disease, termed xe2x80x9cSyndrome Xxe2x80x9d or xe2x80x9cMetabolic Syndromexe2x80x9d, is manifested by defective glucose metabolism (insulin resistance), elevated blood pressure (hypertension), and a blood lipid imbalance (dyslipidemia). See e.g. Reaven, 1993, Annu. Rev. Med. 44:121-131.\nThe evidence linking elevated serum cholesterol to coronary heart disease is overwhelming. Circulating cholesterol is carried by plasma lipoproteins, which are particles of complex lipid and protein composition that transport lipids in the blood. Low density lipoprotein (LDL) and high density lipoprotein (HDL) are the major cholesterol-carrier proteins. LDL are believed to be responsible for the delivery of cholesterol from the liver, where it is synthesized or obtained from dietary sources, to extrahepatic tissues in the body. The term xe2x80x9creverse cholesterol transportxe2x80x9d describes the transport of cholesterol from extrahepatic tissues to the liver, where it is catabolized and eliminated. It is believed that plasma HDL particles play a major role in the reverse transport process, acting as scavengers of tissue cholesterol. HDL is also responsible for the removal non-cholesterol lipid, oxidized cholesterol and other oxidized products from the bloodstream.\nAtherosclerosis, for example, is a slowly progressive disease characterized by the accumulation of cholesterol within the arterial wall. Compelling evidence supports the belief that lipids deposited in atherosclerotic lesions are derived primarily from plasma apolipoprotein B (apo B)-containing lipoproteins, which include chylomicrons, CLDL, IDL and LDL. The apo B-containing lipoprotein, and in particular LDL, has popularly become known as the xe2x80x9cbadxe2x80x9d cholesterol. In contrast, HDL serum levels correlate inversely with coronary heart disease. Indeed, high serum levels of HDL is regarded as a negative risk factor. It is hypothesized that high levels of plasma HDL is not only protective against coronary artery disease, but may actually induce regression of atherosclerotic plaque (e.g., see Badimon et al., 1992, Circulation 86:(Suppl. III)86-94; Dansky and Fisher, 1999, Circulation 100:1762-3.). Thus, HDL has popularly become known as the xe2x80x9cgoodxe2x80x9d cholesterol.\nThe fat-transport system can be divided into two pathways: an exogenous one for cholesterol and triglycerides absorbed from the intestine and an endogenous one for cholesterol and triglycerides entering the bloodstream from the liver and other non-hepatic tissue.\nIn the exogenous pathway, dietary fats are packaged into lipoprotein particles called chylomicrons, which enter the bloodstream and deliver their triglycerides to adipose tissue for storage and to muscle for oxidation to supply energy. The remnant of the chylomicron, which contains cholesteryl esters, is removed from the circulation by a specific receptor found only on liver cells. This cholesterol then becomes available again for cellular metabolism or for recycling to extrahepatic tissues as plasma lipoproteins.\nIn the endogenous pathway, the liver secretes a large, very-low-density lipoprotein particle (VLDL) into the bloodstream. The core of VLDL consists mostly of triglycerides synthesized in the liver, with a smaller amount of cholesteryl esters either synthesized in the liver or recycled from chylomicrons. Two predominant proteins are displayed on the surface of VLDL, apolipoprotein B-100 (apo B-100) and apolipoprotein E (apo E), although other apolipoproteins are present, such as apolipoprotein CIII (apo CIII) and apolipoprotein CII (apo CII). When a VLDL reaches the capillaries of adipose tissue or of muscle, its triglyceride is extracted. This results in the formation of a new kind of particle called intermediate-density lipoprotein (IDL) or VLDL remnant, decreased in size and enriched in cholesteryl esters relative to a VLDL, but retaining its two apoproteins.\nIn human beings, about half of the IDL particles are removed from the circulation quickly, generally within two to six hours of their formation. This is because IDL particles bind tightly to liver cells, which extract IDL cholesterol to make new VLDL and bile acids. The IDL not taken up by the liver is catabolized by the hepatic lipase, an enzyme bound to the proteoglycan on liver cells. Apo E dissociates from IDL as it is transformed to LDL. Apo B-100 is the sole protein of LDL.\nPrimarily, the liver takes up and degrades circulating cholesterol to bile acids, which are the end products of cholesterol metabolism. The uptake of cholesterol-containing particles is mediated by LDL receptors, which are present in high concentrations on hepatocytes. The LDL receptor binds both apo E and apo B-100 and is responsible for binding and removing both IDL and LDL from the circulation. IN addition, remnant receptors are responsible for clearing chylomicrons and VLDL remnants i.e., IDL). However, the affinity of apo E for the LDL receptor is greater than that of apo B-100. As a result, the LDL particles have a much longer circulating life span than IDL particles; LDL circulates for an average of two and a half days before binding to the LDL receptors in the liver and other tissues. High serum levels of LDL, the xe2x80x9cbadxe2x80x9d cholesterol, are positively associated with coronary heart disease. For example, in atherosclerosis, cholesterol derived from circulating LDL accumulates in the walls of arteries. This accumulation forms bulky plaques that inhibit the flow of blood until a clot eventually forms, obstructing an artery and causing a heart attack or stroke.\nUltimately, the amount of intracellular cholesterol liberated from the LDL controls cellular cholesterol metabolism. The accumulation of cellular cholesterol derived from VLDL and LDL controls three processes. First, it reduces the cell\"\"s ability to make its own cholesterol by turning off the synthesis of HMGCoA reductase, a key enzyme in the cholesterol biosynthetic pathway. Second, the incoming LDL-derived cholesterol promotes storage of cholesterol by the action of ACAT, the cellular enzyme that converts cholesterol into cholesteryl esters that are deposited in storage droplets. Third, the accumulation of cholesterol within the cell drives a feedback mechanism that inhibits cellular synthesis of new LDL receptors. Cells, therefore, adjust their complement of LDL receptors so that enough cholesterol is brought in to meet their metabolic needs, without overloading (for a review, see Brown and Goldstein, In, The Pharmacological Basis Of Therapeutics, 8th Ed., Goodman and Gilman, Pergaman Press, NY, 1990, Ch. 36, pp. 874-896).\nHigh levels of apo B-containing lipoproteins can be trapped in the subendothelial space of an artery and undergo oxidation. The oxidized lipoprotein is recognized by scavenger receptors on macrophages. Binding of oxidized lipoprotein to the scavenger receptors can enrich the macrophages with cholesterol and cholesteryl esters independently of the LDL receptor. Macrophages can also produce cholesteryl esters by the action of ACAT. LDL can also be complexed to a high molecular weight glycoprotein called apolipoprotein(a), also known as apo(a), through a disulfide bridge. The LDL-apo(a) complex is known as Lipoprotein(a) or Lp(a). Elevated levels of Lp(a) are detrimental, having been associated with atherosclerosis, coronary heart disease, myocardial infarcation, stroke, cerebral infarction, and restenosis following angioplasty.\nPeripheral (non-hepatic) cells predominantly obtain their cholesterol from a combination of local synthesis and uptake of preformed sterol from VLDL and LDL. Cells expressing scavenger receptors, such as macrophages and smooth muscle cells, can also obtain cholesterol from oxidized apo B-containing lipoproteins. In contrast, reverse cholesterol transport (RCT) is the pathway by which peripheral cell cholesterol can be returned to the liver for recycling to extrahepatic tissues, hepatic storage, or excretion into the intestine in bile. The RCT pathway represents the only means of eliminating cholesterol from most extrahepatic tissues and is crucial to maintenance of the structure and function of most cells in the body.\nThe enzyme in blood involved in the RCT pathway, lecithin:cholesterol acyltransferase (LCAT), converts cell-derived cholesterol to cholesteryl esters, which are sequestered in HDL destined for removal. LCAT is produced mainly in the liver and circulates in plasma associated with the HDL fraction. Cholesterol ester transfer protein (CETP) and another lipid transfer protein, phospholipid transfer protein (PLTP), contribute to further remodeling the circulating HDL population (see for example Bruce et al., 1998, Annu. Rev. Nutr. 18:297-330). PLTP supplies lecithin to HDL, and CETP can move cholesteryl ester made by LCAT to other lipoproteins, particularly apoB-containing lipoproteins, such as VLDL. HDL triglyceride can be catabolized by the extracellular hepatic triglyceride lipase, and lipoprotein cholesterol is removed by the liver via several mechanisms.\nEach HDL particle contains at least one molecule, and usually two to four molecules, of apolipoprotein (apo A-I). Apo A-I is synthesized by the liver and small intestine as preproapolipoprotein which is secreted as a proprotein that is rapidly cleaved to generate a mature polypeptide having 243 amino acid residues. Apo A-I consists mainly of a 22 amino acid repeating segment, spaced with helix-breaking proline residues. Apo A-I forms three types of stable structures with lipids: small, lipid-poor complexes referred to as pre-beta-1 HDL; flattened discoidal particles, referred to as pre-beta-2 HDL, which contain only polar lipids (e.g., phospholipid and cholesterol); and spherical particles containing both polar and nonpolar lipids, referred to as spherical or mature HDL (HDL3 and HDL2). Most HDL in the circulating population contains both apo A-I and apo A-II, a second major HDL protein. This apo A-I- and apo A-II-containing fraction is referred to herein as the AI/AII-HDL fraction of HDL. But the fraction of HDL containing only apo A-I, referred to herein as the AI-HDL fraction, appears to be more effective in RCT. Certain epidemiologic studies support the hypothesis that the AI-HDL fraction is antiartherogenic (Parra et al., 1992, Arterioscler. Thromb. 12:701-707; Decossin et al., 1997, Eur. J. Clin. Invest. 27:299-307).\nAlthough the mechanism for cholesterol transfer from the cell surface is unknown, it is believed that the lipid-poor complex, pre-beta-1 HDL, is the preferred acceptor for cholesterol transferred from peripheral tissue involved in RCT. Cholesterol newly transferred to pre-beta-1 HDL from the cell surface rapidly appears in the discoidal pre-beta-2 HDL. PLTP may increase the rate of disc formation (Lagrost et al., 1996, J. Biol. Chem. 271:19058-19065), but data indicating a role for PLTP in RCT is lacking. LCAT reacts preferentially with discoidal and spherical HDL, transferring the 2-acyl group of lecithin or phosphatidylethanolamine to the free hydroxyl residue of fatty alcohols, particularly cholesterol, to generate cholesteryl esters (retained in the HDL) and lysolecithin. The LCAT reaction requires an apoliprotein such apo A-I or apo A-IV as an activator. ApoA-I is one of the natural cofactors for LCAT. The conversion of cholesterol to its HDL-sequestered ester prevents re-entry of cholesterol into the cell, resulting in the ultimate removal of cellular cholesterol. Cholesteryl esters in the mature HDL particles of the AI-HDL fraction are removed by the liver and processed into bile more effectively than those derived from the AI/AII-HDL fraction. This may be due, in part, to the more effective binding of AI-HDL to the hepatocyte membrane. Several HDL receptor receptors have been identified, the most well characterized of which is the scavenger receptor class B, type I (SR-BI) (Acton et al., 1996, Science 271:518-520). The SR-BI is expressed most abundantly in steroidogenic tissues (e.g., the adrenals), and in the liver (Landshulz et al., 1996, J. Clin. Invest. 98:984-995; Rigotti et al., 1996, J. Biol. Chem. 271:33545-33549). Other proposed HDL receptors include HB1 and HB2 (Hidaka and Fidge, 1992, Biochem J. 15:161-7; Kurata et al., 1998, J. Atherosclerosis and Thrombosis 4:112-7).\nWhile there is a consensus that CETP is involved in the metabolism of VLDL- and LDL-derived lipids, its role in RCT remains controversial. However, changes in CETP activity or its acceptors, VLDL and LDL, play a role in xe2x80x9cremodelingxe2x80x9d the HDL population. For example, in the absence of CETP, the HDL becomes enlarged particles that are poorly removed from the circulation (for reviews on RCT and HDLs, see Fielding and Fielding, 1995, J. Lipid Res. 36:211-228; Barrans et al., 1996, Biochem. Biophys. Acta. 1300:73-85; Hirano et al., 1997, Arterioscler. Thromb. Vasc. Biol. 17:1053-1059).\nHDL is not only involved in the reverse transport of cholesterol, but also plays a role in the reverse transport of other lipids, i.e., the transport of lipids from cells, organs, and tissues to the liver for catabolism and excretion. Such lipids include sphingomyelin, oxidized lipids, and lysophophatidylcholine. For example, Robins and Fasulo (1997, J. Clin. Invest. 99:380-384) have shown that HDL stimulates the transport of plant sterol by the liver into bile secretions.\nPeroxisome proliferators are a structurally diverse group of compounds that, when administered to rodents, elicit dramatic increases in the size and number of hepatic and renal peroxisomes, as well as concomitant increases in the capacity of peroxisomes to metabolize fatty acids via increased expression of the enzymes required for the xcex2-oxidation cycle (Lazarow and Fujiki, 1985, Ann. Rev. Cell Biol. 1:489-530; Vamecq and Draye, 1989, Essays Biochem. 24:1115-225; and Nelali et al., 1988, Cancer Res. 48:5316-5324). Chemicals included in this group are the fibrate class of hypolipidermic drugs, herbicides, and phthalate plasticizers (Reddy and Lalwani, 1983, Crit. Rev. Toxicol. 12:1-58). Peroxisome proliferation can also be elicited by dietary or physiological factors, such as a us high-fat diet and cold acclimatization.\nInsight into the mechanism whereby peroxisome proliferators exert their pleiotropic effects was provided by the identification of a member of the nuclear hormone receptor superfamily activated by these chemicals (Isseman and Green, 1990, Nature 347:645-650). This receptor, termed peroxisome proliferator activated receptor a (PPAR60 ), was subsequently shown to be activated by a variety of medium and long-chain fatty acids. PPARxcex1 activates transcription by binding to DNA sequence elements, termed peroxisome proliferator response elements (PPRE), in the form of a heterodimer with the retinoid X receptor (RXR). RXR is activated by 9-cis retinoic acid (see Kliewer et al., 1992, Nature 358:771-774; Gearing et al., 1993, Proc. Natl. Acad. Sci. USA 90:1440-1444, Keller et al., 1993, Proc. Natl. Acad. Sci. USA 90:2160-2164; Heyman et al., 1992, Cell 68:397-406, and Levin et al., 1992, Nature 355:359-361). Since the discovery of PPARxcex1, additional isoforms of PPAR have been identified, e.g., PPARxcex2, PPARxcex3 and PPARxcex4, which are have similar functions and are similarly regulated.\nPPREs have been identified in the enhancers of a number of genes encoding proteins that regulate lipid metabolism. These proteins include the three enzymes required for peroxisomal xcex2-oxidation of fatty acids; apolipoprotein A-I; medium-chain acyl-CoA dehydrogenase, a key enzyme in mitochondrial xcex2-oxidation; and aP2, a lipid binding protein expressed exclusively in adipocytes (reviewed in Keller and Whali, 1993, TEM, 4:291-296; see also Staels and Auwerx, 1998, Atherosclerosis 137 Suppl:S19-23). The nature of the PPAR target genes coupled with the activation of PPARs by fatty acids and hypolipidemic drugs suggests a physiological role for the PPARs in lipid homeostasis.\nPioglitazone, an antidiabetic compound of the thiazolidinedione class, was reported to stimulate expression of a chimeric gene containing the enhancer/promoter of the lipid-binding protein aP2 upstream of the chloroamphenicol acetyl transferase reporter gene (Harris and Kletzien, 1994, Mol. Pharmacol 45:439-445). Deletion analysis led to the identification of an approximately 30 bp region responsible for pioglitazone responsiveness. In an independent study, this 30 bp fragment was shown to contain a PPRE (Tontonoz et al., 1994, Nucleic Acids Res. 22:5628-5634). Taken together, these studies suggested the possibility that the thiazolidinediones modulate gene expression at the transcriptional level through interactions with a PPAR and reinforce the concept of the interrelatedness of glucose and lipid metabolism.\nIn the past two decades or so, the segregation of cholesterolemic compounds into HDL and LDL regulators and recognition of the desirability of decreasing blood levels of the latter has led to the development of a number of drugs. However, many of these drugs have undesirable side effects and/or are contraindicated in certain patients, particularly when administered in combination with other drugs.\nBile-acid-binding resins are a class of drugs that interrupt the recycling of bile acids from the intestine to the liver. Examples of bile-acid-binding resins are cholestyramine (QUESTRAN LIGHT, Bristol-Myers Squibb), and colestipol hydrochloride (COLESTID, Pharmacia and Upjohn Company). When taken orally, these positively charged resins bind to negatively charged bile acids in the intestine. Because the resins cannot be absorbed from the intestine, they are excreted, carrying the bile acids with them. The use of such resins, however, at best only lowers serum cholesterol levels by about 20%. Moreover, their use is associated with gastrointestinal side-effects, including constipation and certain vitamin deficiencies. Moreover, since the resins bind to drugs, other oral medications must be taken at least one hour before or four to six hours subsequent to ingestion of the resin, complicating heart patients\"\" drug regimens.\nThe statins are inhibitors of cholesterol synthesis. Sometimes, the statins are used in combination therapy with bile-acid-binding resins. Lovastatin (MEVACOR, Merck and Co., Inc.), a natural product derived from a strain of Aspergillus; pravastatin (PRAVACHOL, Bristol-Myers Squibb Co.); and atorvastatin (LIPITOR, Warner Lambert) block cholesterol synthesis by inhibiting HMGCoA, the key enzyme involved in the cholesterol biosynthetic pathway. Lovastatin significantly reduces serum cholesterol and LDL-serum levels. It also slows progression of coronary atherosclerosis. However, serum HDL levels are only slightly increased following lovastatin administration. The mechanism of the LDL-lowering effect may involve both reduction of VLDL concentration and induction of cellular expression of LDL-receptor, leading to reduced production and/or increased catabolism of LDL. Side effects, including liver and kidney dysfunction are associated with the use of these drugs.\nNiacin, also known as nicotinic acid, is a water-soluble vitamin B-complex used as a dietary supplement and antihyperlipidemic agent. Niacin diminishes production of VLDL and is effective at lowering LDL. It is used in combination with bile-acid-binding resins. Niacin can increase HDL when administered at therapeutically effective doses; however, its usefulness is limited by serious side effects.\nFibrates are a class of lipid-lowering drugs used to treat various forms of hyperlipidemia, elevated serum triglycerides, which may also be associated with hypercholesterolemia. Fibrates appear to reduce the VLDL fraction and modestly increase HDL; however, the effects of these drugs on serum cholesterol is variable. In the United States, fibrates have been approved for use as antilipidemic drugs, but have not received approval as hypercholesterolemia agents. For example, clofibrate (ATROMID-S, Wyeth-Ayerst Laboratories) is an antilipidemic agent that acts to lower serum triglycerides by reducing the VLDL fraction. Although ATROMID-S may reduce serum cholesterol levels in certain patient subpopulations, the biochemical response to the drug is variable, and is not always possible to predict which patients will obtain favorable results. ATROMID-S has not been shown to be effective for prevention of coronary heart disease. The chemically and pharmacologically related drug, gemfibrozil (LOPID, Parke-Davis), is a lipid regulating agent which moderately decreases serum triglycerides and VLDL cholesterol. LOPID also increases HDL cholesterol, particularly the HDL2 and HDL3 subfractions, as well as both the AI/AII-HDL fraction. However, the lipid response to LOPID is heterogeneous, especially among different patient populations. Moreover, while prevention of coronary heart disease was observed in male patients between the ages of 40 and 55 without history or symptoms of existing coronary heart disease, it is not clear to what extent these findings can be extrapolated to other patient populations (e.g., women, older and younger males). Indeed, no efficacy was observed in patients with established coronary heart disease. Serious side-effects are associated with the use of fibrates, including toxicity; malignancy, particularly malignancy of gastrointestinal cancer; gallbladder disease; and an increased incidence in non-coronary mortality. These drugs are not indicated for the treatment of patients with high LDL or low HDL as their only lipid abnormality.\nOral estrogen replacement therapy may be considered for moderate hypercholesterolemia in post-menopausal women. However, increases in HDL may be accompanied with an increase in triglycerides. Estrogen treatment is, of course, limited to a specific patient population, postmenopausal women, and is associated with serious side effects, including induction of malignant neoplasms; gall bladder disease; thromboembolic disease; hepatic adenoma; elevated blood pressure; glucose intolerance; and hypercalcemia.\nLong chain carboxylic acids, particularly long chain xcex1,xcfx89-dicarboxylic acids with distinctive substitution patterns, and their simple derivatives and salts, have been disclosed for treating atherosclerosis, obesity, and diabetes (See, e.g., Bisgaier et al., 1998, J. Lipid Res. 39:17-30, and references cited therein; International Patent Publication WO 98/30530; U.S. Pat. No. 4,689,344; International Patent Publication WO 99/00116; and U.S. Pat. No. 5,756,344). However, some of these compounds, for example the xcex1,xcfx89-dicarboxylic acids substituted at their xcex1,xcex1xe2x80x2-carbons (U.S. Pat. No. 3,773,946), while having serum triglyceride and serum cholesterol-lowering activities, have no value for treatment of obesity and hypercholesterolemia (U.S. Pat. No. 4,689,344).\nU.S. Pat. No. 4,689,344 discloses xcex2,xcex2,xcex2xe2x80x2,xcex2xe2x80x2-tetrasubstituted-xcex1,xcfx89-alkanedioic acids that are optionally substituted at their xcex1,xcex1,xcex1xe2x80x2,xcex1xe2x80x2 positions, and alleges that they are useful for treating obesity, hyperlipidemia, and diabetes. According to this reference, both triglycerides and cholesterol are lowered significantly by compounds such as 3,3,14,14-tetramethylhexadecane-1,16-dioic acid. U.S. Pat. No. 4,689,344 further discloses that the xcex2,xcex2,xcex2xe2x80x2,xcex2xe2x80x2-tetramethyl-alkanediols of U.S. Pat. No. 3,930,024 also are not useful for treating hypercholesterolemia or obesity.\nOther compounds are disclosed in U.S. Pat. No. 4,711,896. In U.S. Pat. No. 5,756,544, xcex1,xcfx89-dicarboxylic acid-terminated dialkane ethers are disclosed to have activity in lowering certain plasma lipids, including Lp(a), triglycerides, VLDL-cholesterol, and LDL-cholesterol, in animals, and elevating others, such as HDL-cholesterol. The compounds are also stated to increase insulin sensitivity. In U.S. Pat. No. 4,613,593, phosphates of dolichol, a polyprenol isolated from swine liver, are stated to be useful in regenerating liver tissue, and in treating hyperuricuria, hyperlipemia, diabetes, and hepatic diseases in general.\nU.S. Pat. No. 4,287,200 discloses azolidinedione derivatives with anti-diabetic, hypolipidemic, and anti-hypertensive properties. However, these administration of these compounds to patients can produce side effects such as bone marrow depression, and both liver and cardiac cytotoxicity. Further, the compounds disclosed by U.S. Pat. No. 4,287,200 stimulate weight gain in obese patients.\nIt is clear that none of the commercially available cholesterol management drugs has a general utility in regulating lipid, lipoprotein, insulin and glucose levels in the blood. Thus, compounds that have one or more of these utilities are clearly needed. Further, there is a clear need to develop safer drugs that are efficacious at lowering serum cholesterol, increasing HDL serum levels, preventing coronary heart disease, and/or treating existing disease such as atherosclerosis, obesity, diabetes, and other diseases that are affected by lipid metabolism and/or lipid levels. There is also is a clear need to develop drugs that may be used with other lipid-altering treatment regimens in a synergistic manner. There is still a further need to provide useful therapeutic agents whose solubility and Hydrophile/Lipophile Balance (HLB) can be readily varied.\nCitation or identification of any reference in Section 2 of this application is not an admission that such reference is available as prior art to the present invention.\nIn one embodiment, the invention provides novel compounds having the general formula I: \nand pharmaceutically acceptable salts thereof, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nn and m are independent integers ranging from 0 to 4;\nK1 and K2 are independently selected from the group consisting of xe2x80x94CH2OH, xe2x80x94C(O)OH, xe2x80x94CHO, xe2x80x94C(O)OR5, xe2x80x94OC(O)R5, xe2x80x94SO3H, \nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl;\neach R6 is independently selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl;\nR7 is selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl; and\nwith the proviso that when n and m are both 1 or both 0, then K1 and K2 are not both X, wherein X is selected from the group consisting of xe2x80x94COOH, xe2x80x94C(O)OR5, \nIn another embodiment, the invention provides novel compounds having the general formula I, and pharmaceutically acceptable salts thereof, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nn and m are independent integers ranging from 0 to 4;\nK1 and K2 are independently selected from the group consisting of xe2x80x94CH2OH, xe2x80x94C(O)OH, xe2x80x94CHO, xe2x80x94C(O)OR5, xe2x80x94OC(O)R5, xe2x80x94SO3H, \nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl;\neach R6 is independently selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl;\nR7 is selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl; and\nwith the proviso that when n and m are both 1 or both 0, then K1 and K2 are not both X, wherein X is selected from the group consisting of xe2x80x94COOH, xe2x80x94C(O)OR5, \nIn yet another embodiment, the invention provides novel compounds having the general formula I, and pharmaceutically acceptable salts thereof, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nn and m are independent integers ranging from 0 to 4;\nK1 is selected from the group consisting of xe2x80x94CH2OH, xe2x80x94OC(O)R5, xe2x80x94CHO, xe2x80x94SO3H, \nK2 is selected from the group consisting of xe2x80x94CH2OH, xe2x80x94C(O)OH, xe2x80x94CHO, xe2x80x94C(O)OR5, xe2x80x94OC(O)R5, xe2x80x94SO3H, \nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl;\neach R6 is independently selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl;\nR7 is selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl; and\nwith the proviso that when n and m are both 1 or both 0, then K1 and K2 are not both X, wherein X is selected from the group consisting of xe2x80x94COOH, xe2x80x94C(O)OR5, \nIn yet another embodiment, the invention provides novel compounds having the general formula I and pharmaceutically acceptable salts thereof, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nn and m are independent integers ranging from 0 to 4;\nK1 and K2 are independently selected from the group consisting of xe2x80x94CH2OH, xe2x80x94OC(O)R5, xe2x80x94CHO, xe2x80x94SO3H, \nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl;\neach R6 is independently selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl;\nR7 is selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl; and\nwith the proviso that when n and m are both 1 or both 0, then K1 and K2 are not both X, wherein X is selected from the group consisting of xe2x80x94COOH, xe2x80x94C(O)OR5, \nIn still another embodiment, the invention provides novel compounds having the general formula I, and pharmaceutically acceptable salts thereof, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nn and m are independent integers ranging from 0 to 4;\nK1 and K2 are independently xe2x80x94CH2OH or xe2x80x94OC(O)R5; and\nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl.\nThe compounds of formula I and pharmaceutically acceptable salts thereof are useful for treating or preventing cardiovascular diseases, dyslipidemias, dyslipoproteinemias, disorders of glucose metabolism, Alzheimer\"\"s Disease, Syndrome X, PPAR-associated disorders, septicemia, thrombotic disorders, obesity, pancreatitis, hypertension, renal diseases, cancer, inflammation, or impotence.\nIn another embodiment, the invention comprises a compound of the formula IV: \nwherein:\nn is an integer ranging from 1 to 4;\nK1 selected from the group consisting of xe2x80x94CH2OH, xe2x80x94C(O)OH, xe2x80x94CHO, xe2x80x94C(O)OR5, xe2x80x94OC(O)R5, xe2x80x94SO3H, \nR1, and R2 are independently selected from the group consisting of (C1-C6)alkyl, C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl;\neach R6 is independently selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl;\nR7 is selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl; and\nW is selected from the group consisting of H, (C1-C6)alkyl, and a hydroxy protecting group.\nIn another embodiment, the invention provides a compound of the formula V: \nwherein:\nn is an integer ranging from 1 to 4;\nK1 selected from the group consisting of xe2x80x94CH2OH, xe2x80x94C(O)OH, xe2x80x94CHO, xe2x80x94C(O)OR5, xe2x80x94OC(O)R5, xe2x80x94SO3H, \nR3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\nR5 is selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl;\neach R6 is independently selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl;\nR7 is selected from the group consisting of H, (C1-C6)alkyl, (C2-C6)alkenyl, and (C2-C6)alkynyl; and\nHal is selected from the group consisting of chloro, bromo, and iodo.\nThe compounds of formulas IV and V are useful as intermediates for synthesizing the compounds of formula I.\nIn still another embodiment, the invention provides a method for the synthesis of a compound of a formula II: \ncomprising (a) contacting in the presence of a base a compound of a formula XXIV: \nwith a compound of a formula XXVIII: \nto provide a compound of a formula XXIX: \nand (b) deprotecting the compound of the formula XXIX to provide the compound of the formula II, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH; and\nPG is a hydroxy protecting group.\nIn still another embodiment, the invention provides a method for the synthesis of a compound of formula III: \ncomprising contacting a compound of a formula of formula VI: \nwith a reducing agent, wherein:\nR1, R2, R3, and R4 are independently selected from the group consisting of (C1-C6)alkyl, (C2-C6)alkenyl, (C2-C6)alkynyl, phenyl, and benzyl; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group; or R1, R2, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group and R3, R4, and the carbon to which they are attached are taken together to form a (C3-C7)cycloalkyl group, with the proviso that none of R1, R2, R3, or R4 is xe2x80x94(CH2)0-4Cxe2x89xa1CH;\neach R10 is independently selected from the group consisting of xe2x80x94H, xe2x80x94OH, C1-C8)alkoxy, (C6)aryloxy, xe2x80x94O(C2-C6)alkenyl, xe2x80x94Oxe2x80x94(C2-C6)alkynyl, halo; and\nn and m are independent integers ranging from 0 to 4.\nThe present invention further provides compositions comprising a compound of the formula I or a pharmaceutically acceptable salt thereof; and a pharmaceutically acceptable vehicle. These compositions are useful for treating or preventing a disease or disorder selected from the group consisting of a cardiovascular disease, dyslipidemia, dyslipoproteinemia, a disorder of glucose metabolism, Alzheimer\"\"s Disease, Syndrome X, a PPAR-associated disorder, septicemia, a thrombotic disorder, obesity, pancreatitis, hypertension, a renal disease, cancer, inflammation, and impotence. These composition are also useful for reducing the fat content of meat in livestock and reducing the cholesterol content of eggs.\nThe present invention provides a method for treating or preventing a cardiovascular disease, dyslipidemia, dyslipoproteinemia, a disorder of glucose metabolism, Alzheimer\"\"s Disease, Syndrome X, a PPAR-associated disorder, septicemia, a thrombotic disorder, obesity, pancreatitis, hypertension, a renal disease, cancer, inflammation, and impotence, comprising administering to a patient in need of such treatment or prevention a therapeutically effective amount of a composition comprising a compound of formula I, or a pharmaceutically acceptable salt thereof; and a pharmaceutically acceptable vehicle.\nThe present invention further provides a method for reducing the fat content of meat in livestock comprising administering to livestock in need of such fat-content reduction a therapeutically effective amount of a composition comprising a compound of formula I or a pharmaceutically acceptable salt thereof; and a pharmaceutically acceptable vehicle.\nThe present invention provides a method for reducing the cholesterol content of a fowl egg comprising administering to a fowl species a therapeutically effective amount of a compound of formula I or a pharmaceutically acceptable salt thereof; and a pharmaceutically acceptable vehicle.\nThe present invention may be understood more fully by reference to the figures, detailed description, and examples, which are intended to exemplify non-limiting embodiments of the invention."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to call handling by sequenced applications, and, in particular, to a system and method for bypassing a sequenced application if inadequate system resources are available.\n2. Description of Related Art\nSession Initiation Protocol (SIP) is an open signaling protocol for establishing many kinds of real-time communication sessions. Examples of the types of communication sessions that may be established using SIP include voice, video, and/or instant messaging. These communication sessions may be carried out on any type of communication device such as a personal computer, laptop computer, Personal Digital Assistant, telephone, mobile phone, cellular phone, or the like. One key feature of SIP is its ability to use an end-user's Address of Record (AOR) as a single unifying public address for all communications. Thus, in a world of SIP-enhanced communications, a user's AOR becomes their single address that links the user to all of the communication devices associated with the user. Using the AOR, a caller can reach any one of the user's communication devices, also referred to as User Agents (UAs) without having to know each of the unique device addresses or phone numbers.\nMany SIP communications are enhanced by virtue of the fact that an application is inserted or included into the communication session during the establishment of that session. The incorporation of applications into a communication session is typically referred to as application sequencing because the applications are sequentially invoked during the establishment of the communication session. In some instances the applications are owned and operated by an enterprise that is administering the SIP network. In some instances, the applications may be provided by third-party vendors. In either event, the traditional way in which applications were included in the communication session was during the communication session establishment stage so that these applications can insert themselves into the signaling and media path of the communication session.\nExemplary types of applications that may be utilized for a communication session include, without limitation, call recording applications, communication log services, conferencing applications, security applications, encryption applications, collaboration applications, whiteboard applications, mobility applications, presence applications, media applications, messaging applications, bridging applications, and any other type of application that can supplement or enhance communications.\nSession managers such as Avaya Aura™ for enterprise telephony networks allow sequencing of network applications at session origination and termination phases, in order to affect the way the session is routed in the network. During the origination phase, the call is routed through the servers associated with the calling party's profile. During the termination phase, the call is routed through the servers associated with the called party's profile. The selection of network applications that are used during origination and termination sequencing phases is generally based on the identities of calling and called parties that are involved in a communication session. The list of servers (e.g., feature servers) that forms the sequence vector may be administered at user provisioning time. The calling party's provisioning profile provides the originating sequence vector whereas the called party's profile provides the terminating sequence vector. These two vectors combined provide a map of how a session initiation request travels from calling party to called party in the network.\nSequence vectors presently known in the art are static despite status changes that may occur in the applications that make up the sequence vector. For example, if an application that is in the sequence for a call were to fail, the call setup would fail even if the application is not critical for the completion of the call. However, if the session manager could skip over an application in distress or failure, or execute the application at a reduced level of performance, the call may be completed despite the application being in distress or failure. Currently, there is no easy way for applications that are in the sequence vector to inform a session manager about its current signaling and media status in order to influence the sequencing decision that a session manager makes as a session progresses through the network. Therefore, a need exists to provide dynamic adjustment to sequence vectors in order to provide more robust call handling when a noncritical application is in distress or failure."}
{"text": "The heart is a complicated organ. Stimulation by heart current causes myocardial contraction to pump out blood and makes the blood to circulate in the blood vessel. The heart current passes through the atrium and then the ventricle to cause myocardial contraction. The myocardium rhythmically contracts under the action of heart current. However, the advance in age or the change in pathological and physiological characteristics of the heart causes contraction disorder and abnormal rhythm of the heart. One of the methods to treat arrhythmia is septal ablation, that is, making the abnormal heart current loop open through an interventional surgery. To perform ablation, it is necessary to guide the catheter with an ablation electrode to reach the left atrium. This is usually performed by femoral artery puncture, whereby the electrophysiology catheter with an ablation electrode is delivered to the left atrium under X-ray image observation, the electrode is placed against the orifice of the pulmonary vein, the site of abnormal heart current loop that generates arrhythmia is determined by collecting electric signals, and the discharge energy or other energy generated by the ablation electrode brings about lesion degeneration or coagulative necrosis to the tissue near this site, to thereby cut off the abnormal heart current loop. The energy for use in ablation can be radio frequency current, direct current, microwave, ultrasonic wave or laser.\nThe key to electrophysiology ablation is positioning of the electrode, namely to determine the site to be ablated in a quick and precise manner. U.S. Pat. No. 6,837,886 B2 proposes an electrophysiology ablation device, whose positioning method is carried out by an openable meshed flat disk concealed in the catheter head. During the process of vein puncture, the disk is concealed in the catheter head, and after the catheter head is inserted into the orifice of the pulmonary vein, the disk is opened to abut against the orifice of the pulmonary vein. The disk is braided from a metal wire, which has electrically conductive function, and is itself an electrode that generates electric pulses once the power is on to carry out circumferential ablation on the orifice of the pulmonary vein. That invention is mainly defective in the following aspect: due to the difference in sizes and shapes of the orifices of the patients' pulmonary veins, it is very difficult in practical operation to abut the meshed flat disk completely against the orifice of the pulmonary vein, but instead usually only in a slanting or partial manner, so that it is impossible to carry out ideal circumferential ablation of the orifice of the pulmonary vein, and adversely affecting the therapeutic effect."}
{"text": "Fuel Cell\nA fuel cell is an electrochemical device in which the energy of a chemical reaction is converted directly into electricity. Unlike an electric cell or battery, a fuel cell does not run down or require recharging; it operates as long as the fuel and an oxidizer are supplied continuously from outside the cell.\nA fuel cell consists of an anode to which fuel, commonly hydrogen, or methanol, is supplied; and a cathode-to which an oxidant, commonly air or oxygen, is supplied. The two electrodes of a fuel cell are separated by an ionic conductor electrolyte. In the case of a hydrogen-oxygen fuel cell with an alkali metal hydroxide electrolyte, the anode reaction is 2H2+4OH−→4H2O+4e and the cathode reaction is O2+2H2O+4e→4OH−. The electrons generated at the anode move through an external circuit containing the load and pass to the cathode. The OH— ions generated at the cathode are conducted by the electrolyte to the anode, where they form water by combining with hydrogen. The water produced at the anode is removed continuously in order to avoid flooding the cell. Hydrogen-oxygen fuel cells using ion exchange membranes or immobilized phosphoric acid electrolytes found early use in the Gemini and Apollo space programs, respectively.\nHydrogen Fuel Cell\nThe hydrogen fuel cell is an attractive replacement alternative to the internal combustion engine for producing electricity because it is both highly efficient and weakly polluting. The hydrogen fuel cell has existed for a long time and could produce an estimated 40% of the total energy demand within the next fifty years.\nA conventional hydrogen fuel cell identified with numeral reference 10 is shown in FIG. 1. As shown, a hydrogen fuel cell mainly consists of two electrodes, an anode (2) and a cathode (4), separated by an electrolyte (6). At the anode (2), gaseous hydrogen liberates an electron and a positively charged ion. The electron transits through an external circuit (8) producing an electrical current while the ion diffuses through the electrolyte. At the cathode (4), the electrons combine with the hydrogen ions and oxygen to form water a non-polluting waste product that can be reused. A catalyst is used in order to accelerate this oxidizing reaction.\nHydrogen fuel cells are classified according to operating temperature and electrolyte type. There are five different categories of hydrogen fuel cells according to the type of electrolyte used: phosphoric acid; molten carbonate; solid oxide; proton-exchange membrane; and alkaline.\nIn order to implement hydrogen fuel cells in a day-to-day life, two technological obstacles must be overcome. The first concerns the cell itself and the second concerns the production, storage and transport of the hydrogen that is the energy vector of the process. Up to now, the development of the hydrogen fuel cell has leaped ahead while the hydrogen storage technology lags behind. This technological delay could very well be a determining factor in the marketability of hydrogen fuel cells.\nAlthough hydrogen is a relatively widely present element in nature, it rarely presents itself as a readily available molecularly free gas. Hydrogen is therefore not considered as a source of energy per se, but rather as an energy vector. Various methods are available in order to produce hydrogen. Among these, there is the electrolysis of water or pyrol. However, the only profitable method at the present time is the production of hydrogen from fossil fuels.\nOnce produced, hydrogen may be stored as a liquid at extremely low temperatures or alternatively, under extremely high pressure. However, the high cost of this storing approach favours the use of a technology where the hydrogen is produced in place and at a production rate governed by its consumption rate. Hence, instead of storing liquid hydrogen, a liquid fuel containing hydrogen, for instance liquid hydrocarbons, is transported in place to be transformed in a later stage, and as needed, into hydrogen.\nThe presently available hydrogen fuel cell technology typically uses a liquid fuel containing hydrogen, for example diesel, gasoline, methanol, natural gas, etc. Such liquid fuel is converted by a reformer into a gaseous mixture of hydrogen and carbon dioxide. Then, the hydrogen is extracted from the gaseous mixture and supplied to the hydrogen fuel cell.\nMethanol, for example, which is in liquid form at room temperature, contains a large quantity of hydrogen that may be chemically extracted. The process used to convert hydrocarbons such as methanol into hydrogen and carbon is internal or external reforming. Partial oxidizing and autothermic reforming are other physico-chemical processes that allow for separation of hydrogen from hydrocarbons. Vapour reforming of hydrocarbons (crude or refined oil, natural gas or gasified coal) yields hydrogen-rich mixtures that must be treated further in order to remove carbon monoxide and carbon dioxide which hinder the operation of certain hydrogen fuel cells.\nThe vapour reforming of natural gas involves the catalysed endothermic conversion of light hydrocarbons (including methane through gasoline) with water vapour. The industrial process usually occurs at 850° C., under 2.5 mPa of pressure, according to the following equation:CnHm+nH2OnCO+(n+m/2)H2  (I)                where n is an integer equal or greater than 1;m=2n+2.        \nThe catalyzed exothermic conversion (shift reaction) of the carbon monoxide by-product produces hydrogen according to the following equation:CO+H2OCO2+H2  (II)\nThe CO2 is separated from the gas mixture by absorption processes or membrane separation, and the hydrogen, before being transferred to the hydrogen fuel cell, is treated a second time in order to remove unwanted compounds. Presently, and despite the important costs linked to its extraction, the carbon dioxide is generally released into the atmosphere, thereby contributing to the increase of the greenhouse effect.\nThe partial oxidizing of heavy hydrocarbons involves the exothermic or autothermic conversion of heavy hydrocarbons in the presence of oxygen and water vapour. The following equation describes the phenomenon:CH1,4+0,3H2O+0,4O20,9CO+0,1CO2+H2  (III)\nThe partial oxidizing of coal, except for the initial preparation of coal, follows the same process as the one used for the gasification of heavy hydrocarbons. By-products also include CO2.\nAlthough hydrogen fuel cell technologies presently pollute less than internal combustion engines, these systems when combined with hydrocarbon reforming for example, still produce CO2, the greatest contributor in terms of volume to the greenhouse effect. The growing concerns over climatic changes add to the pressures aimed at reducing such emissions which contribute to the greenhouse effect. The aforementioned need for decarbonization of hydrocarbons during reforming is but one example that warrants the use of a carbon dioxide separation and sequestration technology in order to avoid CO2 dumping into the environment. Also, the presence of all the non-hydrogenated contaminants in the fuel decreases the hydrogen fuel cell's power and efficiency.\nThere already exists, in the prior art, a large number of processes for the production of electricity which are based on the hydrogen fuel cell concept. Among these, many rely on a typically less polluting process than the internal combustion engine but still dump gaseous CO2 back into the atmosphere (JP2001-229946; JP11-229167; WO 01/23302; U.S. Pat. Nos. 6,299,744; 5,746,985) or into the water (U.S. Pat. No. 5,141,823).\nLiterature abounds with a substantial number of systems integrating a treatment intended to reduce CO2 emissions by these types of cells. The United States government already uses this technology for certain space applications and intends of providing such a system for a future Mars mission (U.S. Pat. No. 5,005,787). The carbon dioxide produced by the crew and the hydrogen fuel cell would be introduced into a greenhouse and in return, the oxygen-charged air would be recirculated back into the living quarters and cabins. Japanese documents JP2000-287546, JP11-275965, and JP06-333589 also propose a similar solution in order to definitely remedy the CO2 emission problem stemming from hydrogen fuel cell operation.\nCO2 recycling is also proposed in other applications, often depending on the type of cell used. CA 2,326,024 and JP2001-219271, for example, suggest the recycling of CO2 and its use as a shielding gas during an arc welding process. The molten carbonate cells are a particular case in which CO2 is used in order to enrich the oxidant at the cathode. A number of these technologies aim to convert CO into CO2 first (JP2000-340247; JP2000-302405; JP2000-251918; JP11-199202; TW 432,741), and then to recycle this CO2 toward the hydrogen fuel cell's cathode (JP2000-331698; JP2000-260446; JP11-162488; JP11-067251; JP06-275291; JP02-301968; CA 2,002,227/Absorption).\nSome designers wager on a slightly different know-how, namely the methanation of CO2 (JP04-190831; JP08-069808; WO 01/04045; U.S. Pat. Nos. 5,997,594; 6,221,117). Thus reconversion of carbon dioxide into CH4 permits the recycling of carbon and hydrogen by conventional hydrogen fuel cells. U.S. Pat. No. 6,187,465 puts forward a means for recirculating CO2 for a hydrogen fuel cell without any unwanted emissions.\nAdsorption is a <<classic>> which invariably sustains interest. Zeolites (JP08-069808; EP 0 700 107) occupy a substantial part of this consideration, but other systems also attract attention such as JP09-204925 which uses an ionic resin. The separation of gaseous hydrogen by way of a porous membrane or surface (JP2001-139304/zirconate; CA 2,322,871) allows for the extraction of CO2 which may then be stocked (WO 01/83364), recycled or discarded (U.S. Pat. No. 5,759,712). There is also the PSA (Pressure Swing Adsorption) described for example in JP10-027621; JP62-274561; and U.S. Pat. No. 6,299,994.\nObviously, the separation and sequestration of CO2 remain a largely explored area covering a realm of possibilities. In view of the fact that carbon dioxide reduces the performance of hydrogen fuel cells, numerous processes include either a CO2 <<separation unit>> (JP62-170171; U.S. Pat. No. 5,248,567; CA 2,332,079), a CO2 <<separation apparatus>> (JP06-044993; JP07-169482; JP04-337252; WO 00/03126), or what is known as a <<scrubbing device>> (WO 01/25140; EP 0 677 883; U.S. Pat. No. 4,537,839; JP59-184468). Finally, JP11-162495 refers to a <<decarbonation system>>.\nAbsorption of exhaust fumes is another method which sustains keen interest among many inventors as a number of patents/patent applications demonstrate (JP02-206689; JP60-035470; JP60-241673; JP2001-106503/aqueous sodium hydroxide; JP2001-023677/water or alkaline solution; EP 0 376 219) (JP2000-297656; JP04-190831; JP04-051464; WO 01/04045; WO 01/17896; WO 01/23302; WO 01/65621; CA 1,291,787). Although different means have been employed, the cornerstone of innovation in this field remains the improvement of the solubilization of the gas in question.\nAlso known in the prior art, there are processes aiming at reducing CO2 emissions in general. Examples of those processes are given in U.S. Pat. Nos. 5,514,378; 6,143,556; CA 2,222,030; 6,258,335; and EP 0 991 462.\nU.S. Pat. No. 5,514,378 discloses a process for the reduction of CO2 emissions in an enzymatic photobioreactor using carbonic anhydrase. Carbonic anhydrase is a very reactive enzyme that is common to most animal and plant species and hence, readily available. Trachtenberg (U.S. Pat. No. 6,143,556; CA 2,222,030) describes a system for the treatment of exhaust fumes with an enzyme that may be carbonic anhydrase, but does not suggest any specific application with respect to hydrogen fuel cells. The University of Michigan presents a photobioreactor for the treatment of CO2 with carbonic anhydrase for a medical application as an artificial lung (WO 92/00380; U.S. Pat. No. 5,614,378).\nU.S. Pat. No. 6,258,335 discloses a catalytic process for the removal of exhaust CO2 through chemical fixation. This process comprises the optional use of carbonic anhydrase without any particular consideration for reforming exhaust. Materials Technology Corporation (U.S. Pat. No. 4,743,545) presents a bioreactor with a catalyst, that may be carbonic anhydrase, included inside of a hollow bead.\nEP 0 991 462 in the name of the applicant discloses a counter-current-packed-column for CO2 treatment. In this process, carbonic anhydrase is used in free or immobilized states.\nU.S. Pat. No. 6,110,370 discloses a process where carbonic anhydrase is used to treat water and U.S. Pat. Nos. 4,602,987; 4,761,209; CA 1,294,758 in the name of Aquanautics Technology Corporation disclose a method for extracting and using oxygen from fluids. This system includes a step for the separation of CO2 by using the carbonic anhydrase enzyme and an electrochemical cell.\nAlthough a lot of efforts have been put forward in the development and improvement of hydrogen fuel cell in terms of production and storage of the hydrogen fuel, and in terms of reduction of CO2 emissions, there is still place for improvements in this field. More particularly, there is still a need for a process or an apparatus that would provide a simple and affordable way of producing the required fuel, hydrogen; and that without further polluting the atmosphere with CO2 emissions."}
{"text": "Trench excavation for purposes of installing an underground utility infrastructure is an essential and common activity of any ground development work. The common applications are for various kinds of piping, communications, energy etc. There are different types of digging equipment, and each has its advantages and disadvantages. The common denominator with all equipment manufacturers of mechanical digging equipment is the search for solutions that enable more efficient digging and reduction of the limitations connected with ground hardness.\nThere are two types of equipment in use today, excavators and trenchers. Excavators have a jointed arm powered hydraulically, with a digging tool at one end. The advantage of these machines lies in the fact that they are universal and flexible in the type of digging they can achieve. However, they are not efficient as they do not operate continuousiy, and require three movements, a first movement for digging, a second for removing the soil from the trench and dumping it along the sides of the trench, and a third for returning to the trench for digging. Thus, the effective digging time is approximately only 30% of working time.\nAdditionally, these machines do not move while digging. At the end of a digging cycle, depending on the length of the arm, they must be re-positioned. For this reason, the total effective digging time is reduced below 30% of the working time.\nAn additional limitation of these machines is the ground cutting speed of the digging tool. The speed is very low and thus the digging tool requires a large force to split the ground, and must be very heavy. For these reasons, these machines have difficulty digging into hard and rocky ground.\nIt is possible to equip excavators with a hydraulic hammer to break the rocks, after which another pass is required to clean the trench. Clearly, these extra operations greatly increase the excavation cost.\nTrenchers operate in continuous fashion, thus saving valuable working time. The major disadvantages of these machines are that the digging systems are assembled from a large number of moving parts and chains including a separate soil removal system. Multiplication of systems and moving parts reduces the efficiency and technical reliability, while increasing the price. Further, since trenchers have a large contact area with the ground, and cut on a diagonal, they require large amounts of power, force and weight for digging operations.\nOther types of trenchers include those having a large digging wheel with cutting teeth. These machines are limited in their cutting depth which is always less than the wheel radius. This type of trencher also has a large contact area with the ground, and since the wheel is driven from its center, large amounts of power and large moments of force are required for digging operations. An additional soil removal system is required.\nExamples of patents disclosing trench excavation equipment include Japanese Patent No. 60-25-129 to Miwa, in which a drum-shaped cutter is pivotally supported on the front of a traveler with a conveyer to carry excavated soil to the rear.\nRussian Patent 457777 to Kudra discloses a trench excavator device having inclined knives attached to a screw conveyor, facing the spiral direction to improve operation.\nFrench Patent 2,566,024 to Corneille discloses a narrow trench digger with vertical rotating auger, with a parallel partition behind, equipped with soil loosening tools, such as vibrating vertical toothed bars, teeth on the auger spirals, or vertical rotary cutters.\nBelgian Patent 902104 to Durieux discloses a trench cutting machine having a rotating tool comprising a pipe on which a flat is spiral wound or threaded with cutting or abrasive pieces. The trench cutting machine propels itself along the ground via a winch. Belgian Patent 1005788 to Durieux discloses a rotary trench cutter with a detachable vertical partition, to produce trenches of different sizes.\nTherefore, it would be desirable to provide a trencher for excavating various size trenches, which operates at a high speed with increased capability for use on hard and rocky ground."}
{"text": "High-k materials (k>5.0) are used as insulating spacer materials in semiconductor devices. It is desirable to modify the k-value to improve device performance. In addition, device performance suffers when dopants diffuse into the channel region from the source/drain regions of a semiconductor device. It is desirable to prevent dopant diffusion into the channel region."}
{"text": "A client device may access remotely an application running on a remote server. However, a remote application may be designed for a desktop environment for a large screen, which can make it very difficult for a user at a client device having a small screen to interact with the remote application.\nInstead of remotely accessing a desktop application from a client device, the application may be rewritten to run on the client device. However, client devices are typically slower and have limited resources compared with servers (e.g., desktop computer). As a result, rewriting an application on a server to run on a client device can be a huge undertaking requiring a large investment and long development time. For example, it can be very difficult to port existing code from an operating system of a server to a client device (e.g., smartphone). In many cases it is not practical to rewrite an application from scratch to match the native functionality and user interface of a client device.\nAccordingly, there is a need for systems and methods that facilitate user interaction with an application running on a server from a client device and provide an improved experience for users at the client device."}
{"text": "Ion implantation is a technology widely used in semiconductor device manufacturing. In an ion implantation apparatus, an ion beam is generated and directed at a semiconductor wafer to implant ions therein, thereby forming ion implanted regions having desired conductivity. For different types of ion implanted regions, different ion beam properties, such as beam current and/or beam energy, are used. Various types of ion implantation apparatuses are developed to provide ion beams of different properties. Examples of such ion implantation apparatuses include medium current implanters (MCI), high current implanters (HCI), and high energy implanters (HEI). Ion implantation apparatuses of different types, e.g., MCI and HCI, are often provided as separate tools."}
{"text": "Power switching devices are routinely used in a number of applications; commonly, a power switching device is implemented by means of a transistor—for example, an Insulated Gate Bipolar Transistor (IGBT)—capable of sustaining high voltages (such as up to 100-1,000V) and of driving large currents (such as up to 0.1-10 A). A typical example of an application of the IGBT is in the automotive field, wherein the IGBT can be used to control the ignition sparks of plugs in an internal-combustion engine.\nParticularly, in such application the IGBT is coupled with a primary winding of a transformer; the transformer has a plurality of secondary windings, each one coupled with a respective spark plug. The IGBT is firstly turned on by applying a suitable voltage to its gate terminal. As a result, the IGBT passes from an off (blocking) state—wherein a collector-emitter voltage thereof is about equal to a voltage provided by an automotive battery (typically 12V, with respect to a reference or ground voltage)—to an on state—wherein the same collector-emitter voltage reaches a saturation voltage (such as lower than 1V). In this way, a voltage across the primary winding passes from the ground voltage (i.e., 0V) to approximately the battery voltage (i.e., 12V). This causes the charging of the primary winding with a current having a linear-like pattern. The peak value of this charging current is determined by the length of a time interval, during which the IGBT is kept on. At this point, when the shooting of an ignition spark is required the IGBT is turned off so as to cause an abrupt cut of the charging current. Consequently, an extra-voltage appears across the primary winding; this generates a very high voltage at each secondary winding (of the order of some thousands of volts), which high voltage causes the generation of the ignition spark.\nWhen the IGBT turns on in order to charge the primary winding, the voltage across the primary winding undergoes a sharp variation, having a duration corresponding to a turn-on transient period of the IGBT (while switching from the off state to the on state). Typically, by applying a step voltage to the gate terminal of the IGBT, the duration of the turn-on transient period is of the order of hundreds of nanoseconds. However, this results in a very high incremental ratio ΔV/Δt of the voltage across the primary winding, which generates an overshoot that may cause an undesired ignition spark.\nIn order to solve this problem, the IGBT is generally controlled to obtain a so-called soft turn-on thereof, wherein the collector-emitter voltage of the IGBT is gradually decreased (from 12V to 1V). For this purpose, it is possible to apply a direct turn-on current to the gate terminal of the IGBT; the turn-on current charges corresponding stray capacitors, so as to increase the gate voltage relatively slowly until the IGBT turns on. In this way, the incremental ratio ΔV/Δt of the voltage across the primary winding of the transformer is greatly reduced (thereby avoiding any undesired ignition sparks).\nUnfortunately, the above described soft turn-on procedure may increase a turn-on delay between the application of the signal required to turn-on the IGBT and its actual switching. The turn-on delay may cause a corresponding reduction of the maximum charging current that is reached when the IBGT is turned off to generate the ignition spark, and consequently a reduction of the energy stored in the transformer (for the same time available); this may cause a poor ignition spark when the turn-on delay is too long.\nA solution known in the art for reducing the length of the turn-on delay (without causing any overshoots) consists of pre-charging the gate terminal of the IGBT—to a pre-charging voltage lower than its threshold voltage—before the application of the above-described turn-on current. However, the threshold voltage of the IGBT is strictly related to manufacturing process spreads and to temperature variations. Therefore, the pre-charging voltage takes a worst-case value sufficiently low to ensure that the IGBT is kept off (before applying the turn-on current) in any condition. Accordingly, such solution may not be completely satisfactory, since in most practical situations the length of the turn-on delay remains significantly high.\nIn any case, the length of the turn-on delay varies according to the actual operative conditions. Therefore, it may be impossible to control the generation of the ignition sparks accurately."}
{"text": "Periodic signals are used in a variety of electronic devices. One type of periodic signal are clock signals that can be used to establish the timing of a signal or the timing at which an operation is performed on a signal. For example, data signals are typically coupled to and from memory, such as synchronous dynamic random access memory (“SDRAM”), in synchronism with a clock or data strobe signal. More specifically, read data signals are typically coupled from a memory in synchronism with a read data strobe signal. The read data strobe signal typically has the same phase as the read data signals, and it is normally generated by the same memory device that is outputting the read data signals. Write data signals are typically latched into a memory device in synchronism with a write data strobe signal. The write data strobe signal typically has a phase relative to the write data signals so that a write data strobe signal transitions during a “data eye” occurring at the center of the period in which the write data signals are valid.\nInternal clock signals generated in electronic devices, for example, memory devices or memory controllers, are often synchronized or have some other controlled phase relationships relative to external or internal clock signals. For example, with reference to a memory, a clock signal used for both latching write data and outputting read data may be generated in the memory to which the data are being written. The clock signal is typically generated in the memory device from an internal clock signal that is also derived from the system clock signal.\nVarious techniques can be used to generate a clock signals or read/write data strobe signal. FIG. 1 illustrates a conventional clock circuit 100 providing an output clock signal ICLK to a clock tree circuit 140. When enabled by an active CkEn signal, the clock tree circuit 140 distributes the ICLK signal as a DCLK signal to various circuitry that operate according to the DCLK signal. In FIG. 1, the DCLK signal is provided to data output circuitry 150. In particular, the DCLK signal clocks a data register 154 which provides data to an output buffer 158 to generate a data output signal DQ. The clock circuit 100 generates an ICLK signal that when delayed through the clock tree circuit 140 results in a DCLK signal that is synchronized with a reference clock signal RCLK (and its complement RCLK/). The clock generator includes a delay-locked loop (DLL) 102 and a duty cycle correction (DCC) circuit and output buffer 116. The DLL includes an input buffer 104 that provides a buffered reference clock CLKS to a DLL delay line 108. The delayed buffered reference clock signal is output to the DCC and output buffer 116 for correction of the duty cycle and buffering before being output as the ICLK signal. The ICLK signal is also provided to a model delay 120. The model delay 120 models propagation delay through the output buffer and the clock tree circuit 140. A feedback clock signal FBCLK is output from the model delay 120 and provided through model delay 124 to a phase detector circuit 128 as the DLLFB signal. The model delay 124 models the propagation delay of the input buffer 104. The phase detector circuit 128 detects a phase difference between the CLKS and DLLFB signals. A phase difference signal indicative of the phase difference between the CLKS and DLLFB signals is provided to shift logic 132 that generates a control signal based on the phase difference signal to adjust the delay of the DLL delay line 108. The delay is increased or decreased in order to synchronize the CLKS and DLLFB signals. When synchronized, the clock circuit 100 is said to be “locked.”\nAlthough the timing of the ICLK signal relative to the RCLK (and RCLK/signal) is set by the clock circuit 100 so that the DQ signal output by the clock tree circuit 140 is synchronized with the RCLK signal, there may be a “jitter” in the DQ signal. The clock jitter may be caused by the resolution of adjustment in the DLL delay line 108, and also the resolution of the phase detection by the phase detector circuit 128. Clock jitter may also be caused by varying operating conditions, such as varying power, voltage, and temperature. For example, the ICLK signal may need to be driven over a relatively long signal line to be input in the clock tree circuit 140. Although the propagation delay to the clock tree circuit 140 and through the clock tree circuit 140 is modeled by the feedback model delay 120, the actual propagation delay due to the signal line may vary under changing operating conditions, thus, resulting in clock jitter."}
{"text": "Stack molding is well known in the injection molding art and provides various advantages. In particular, stack molding enables the output of an injection molding machine to be at least doubled without significantly increasing its size or clamping tonnage. Stack molds are typically double or quadruple-level, although there could be any number of stacks in a molding machine. For example, some rubber molds use up to ten levels.\nA double level stack mold generally comprises a stationary first platen, a movable center platen and a movable second platen, with two single face mold plates mounted back to back. A first mold (single cavity or multi-cavity) is defined by one of a mold cavity or core plate which is located on the face of the movable center platen adjacent the stationary first platen and the other of the mold cavity or core plate which is located on the stationary first platen. A second mold is defined by one of a mold cavity or core plate which is located on the other face of the moveable center platen adjacent to the movable second platen and the other of the mold cavity or core plate located on the moveable second platen. The molds are opened and closed by a single machine force actuator (generally a hydraulic ram) applied to the moveable second platen and transferred from the second platen to the center platen by a suitable linkage. In a quadruple stack mold, an additional two moveable platens are provided and mold cavity plates and/or mold core plates are located thereon to define additional molds.\nTo supply molten resin to the cavities of the closed molds, conventional stack molds employ a sprue bar which runs from the machine through the stationary platen to the center platen and which serves as a direct channel between the extruder nozzle of the injection molding machine and the mold's hot runner distributor, which is mounted in the center platen of the stack mold. Alternatively, a movable sprue bar located outside of the mold stack can convey the resin to the center section, as described by Bertschi in U.S. Pat. No. 5,011,646. Sprue bars generally include heaters along their length to maintain the molten state of the resin traveling therethrough and must cope with the relatively high pressure at which the molten resin passes through them.\nFor injection molding applications where there are more than two levels in the stack mold, multiple sprue bars can be used for delivering a split stream of molten resin to the hot runner distributors in the multi-level injection mold. In this case, after the resin stream is split, the sprue bars carry the resin to the hot runner distributors in their respective mold sections comprising the injection mold. With multi-sprue bar applications, a single source injection unit channel is typically used with a machine nozzle that divides the single source channel into a plurality of channels aligned with the individual sprue bars, as described in U.S. Pat. No. 5,522,720 to one of the present inventors, and assigned to the assignee of the present invention.\nIn such cases, the sprue bars are normally attached to the respective mold section to which the molding resin is being delivered. Because injection mold sections in a multi-level stack mold generally move in the longitudinal or vertical direction when the mold is open and closed, the sprue bars must be displaced with the mold sections. Accordingly, the sprue bars are not rigidly attached to their source of resin, i.e.--the machine nozzle or the channel splitting device. Consequently, the sprue bar arrangement must be designed so that the sprue bars will return to their sources of resin and reform a seal therewith at the beginning of each molding cycle.\nIn particular, several design problems are typical for stack molds with more than two levels where resin must flow from a single source injection unit to multiple levels spaced progressively farther from the stationery platen. For example, in a four level stack mold, a sprue bar will feed the first and second level via channels in the mold plate between the two levels and a second sprue bar will feed the third and fourth levels via channels in the mold plate between these two levels. It is desired that sprue bars be as short as possible to reduce pressure losses and to minimize the manufacturing expense of the sprue bars. A further difficulty occurs as, due to the progressive arrangement, the two sprue bars will necessarily be of different lengths and thus the pressure drop that occurs between the inlet end of the sprue bar adjacent to the injection nozzle and the outlet end of the bar is much larger in the longer sprue bar than in the shorter.\nWhen molding shallow parts, and thus opening the mold to a relatively small degree, the length of the sprue bars is generally not large, and the difference in the sprue bar length is relatively small. Consequently, the pressure drop is of minor consequence, generally on the order of 3 to 5 MPa. However, a relatively large pressure drop, on the order of 25 MPa, can occur when molding tall parts because the sprue bars are necessarily longer. This large pressure drop must be compensated for at the injection molding machine and, more importantly, the differential in the pressure drop between the sprue bars can cause insufficient mold packing in the molds furthermost from the injection nozzle.\nAnother difficulty with sprue bars is that variations in their length occur due to thermal expansion effects, as the sprue bars are heated to allow resin to flow through them. Accordingly, when the mold is closed, the position of the end of each sprue bar relative to the stationery platen and the channel splitting nozzle on the injection molding machine will vary, due to these thermal variations and due to variations in the position in which the mold plates close at the various levels. The combination of these variables makes it very difficult to predict the location of the two sprue bar ends each time the mold is closed and the sprue bar is returned to the channel splitting nozzle. Therefore, some resin leakage from the joint between the nozzle and the sprue bars is inevitable. Resin also tends to leak or \"drool\" from the nozzle gates or the open channel of the sprue bars when the mold is opened. This drool cannot be tolerated at any of the parting lines of the mold cavity and core sections. At best, such drool prevents complete mold closing and allows flashing to occur and, at worst, can cause permanent damage that requires expensive repairs.\nU.S. Pat. No. 5,522,720 to one of the present inventors, and assigned to the assignee of the present invention, discloses a nozzle that tolerates misalignment with the two sprue bars while still forming a tight and repeatable seal between the sprue and the nozzle. Although this design overcomes the problem of drool between the nozzle and sprue bar during injection (mold closed), it does not solve the problem of drool from the gates when the mold opens, or the substantial pressure drops to the lengths of the two sprue bars.\nU.S. Pat. No. 4,207,051 to one of the present inventors and assigned to the assignee of the present invention shows a stack mold wherein molten resin is supplied to a center platen through a telescoping tube assembly which is mounted externally to the mold. Essentially, the two tubes form an expandable single sprue bar to deliver molten resin to the mold hot runner and the sprue bar thus need not be detached from the injection nozzle. However, it has proven difficult to construct and operate such a telescoping tube system to accommodate the very high injection pressures (exceeding 20,000 psi) experienced at the nozzle in a multi-cavity, multi-level stack arrangement.\nU.S. Pat. No. 5,458,843 discloses a four level stack mold that utilizes a single sprue bar with feed connectors extending through mold components. Drool is reduced via a valveless anti-drool arrangement whereby a spring-activated extension of an outwardly tapered piston into the manifold flow passage reduces its internal pressure and thereby minimizes backflow and resin drool from the feed connector. However, no provision is made for the possibility of drool in the central distributor side of the feed connector. Therefore, the risk of leakage at the mold parting line still exists.\nU.S. Pat. No. 4,212,626 to Gellert dispenses with sprue bars entirely and instead uses a combination of control valve units abutted together to transfer the pressurized melt through mechanically operated valve gates from the stationary platen, where the machine nozzle resides, to the moving platen, where the hot runner manifold resides. Several problems are inherent in this approach. First, drool at the parting lines is likely to occur over time, as the valve gates must remain aligned at the parting mold faces with extreme precision over millions of injection cycles. Second, melt channel capacity is limited by the size of the valve gates through which the melt must pass through. Therefore, large parts cannot be successfully molded using this arrangement. Third, the mold shut height is much greater to accommodate the arrangement of the valve gate construction, leading to slower cycle times and greater material expense (the platens are thicker and therefore have greater mass).\nU.S. Pat. No. 4,611,983 to Bielfeldt discloses a transfer molding system for fiber-reinforced thermoset resins whereby the molten resin is transferred to an injection cylinder via a feed bore. The injection piston is connected to a telescoping sleeve, so that as the piston moves up inside the injection cylinder to fill the mold cavity with resin, the sleeve also rises and seals off the feed bore. Also, the inner diameter of the sleeve is larger than the root diameter of the injection piston, so that any resin drool flows out of the annular clearance. However, this technique suffers from various disadvantages and does not work if applied to a high injection pressure, multi-cavity, multi-level stack mold arrangement with at least two hot runner systems.\nU.S. Pat. No. 4,586,887 to Gellert shows two opposed hot tip bushings meeting at a stack mold parting line. Each hot tip bushing has a first, inner heating element and a second, outer heating element, each running along the length of the hot tip. A temperature difference is created by heating the inner heating element of the first nozzle and the outer heating element of the second nozzle, which, when combined with the taper of the gates, is intended to reduce tip drool, as the tips freeze off each time the mold faces separate. Excess melt is intended to be absorbed into the system when the mold is opened. However, this arrangement has been less than satisfactory in eliminating tip drool.\nU.S. Pat. No. 4,891,001 to Gellert was directed to overcome some of the drooling problems of U.S. Pat. No. 4,586,887 discussed above. This reference teaches an arrangement wherein, rather than heating the tip with two coaxial heaters running the length of the tips, localized heat control is provided on the tip directly at the orifice in the mold face. The first tip has a heating element proximate the mold face, while the second tip had a heating element that is distal the mold face. A similar principal of creating unequal temperatures of the tips at mold separation is intended to result in freezing. However, this arrangement has also been found to experience unacceptable drool.\nAccordingly, it is desired to have a stack mold and a sprue bar assembly therefor which does not suffer from the above-mentioned or other disadvantages."}
{"text": "1. Field of the Invention\nThe present invention relates to a thin-film patterning method for a magnetoresistive device.\n2. Related Background Art\nConventional thin-film patterning methods for thin-film magnetic heads and the like having magnetoresistive devices, for example, are described in Japanese Patent Application Laid-Open Nos. 2002-175606, 2003-17353, 2003-512941, 2000-76618, 2001-110663, and 5-342527, for example. These methods can perform thin-film patterning of thin-film magnetic heads and the like."}
{"text": "Presently available are devices which generate short laser pulses including a laser oscillator which generates a first laser pulse, an element to compress the pulse by reflecting the first laser pulse by stimulated Brillouin scattering into a laser pulse of a shorter temporal duration, and a primary ray path that includes the laser oscillator and the pulse compressing element and along which the first laser pulse travels from the laser oscillator to the pulse compressing element.\nSuch devices, known from prior art, normally make use of laser oscillators generating the first laser pulse with the temporal duration in the range of nanoseconds, most preferably in the range between approximately a few and and some tens of nanoseconds. Laser pulses generated in that way are shortened by one reflection cycle inside a pulse compressing element.\nThe shortcomings of these devices are that the temporal duration of the shortened or compressed pulse depends on the refractive medium inside the pulse compressing element and--in particular--the fact that there is no possibility to generate laser pulses of a predefined time duration.\nTherefore, this invention has the underlying task to create a device and a process that opens up the possibility of generating laser pulses with a predefined time duration."}
{"text": "1. Field of the Invention\nThe present invention relates generally to ferromagnetic memory and more specifically to ferromagnetic memory utilizing giant magnetoresistance and spin polarization.\n2. Description of the Background Art\nFor many years, random access memory for computers was constructed from magnetic elements. This memory had the advantage of very high reliability, nonvolatility in the event of power loss and infinite lifetime under use. Since this memory was hand assembled from three-dimensional ferrite elements, it was eventually supplanted by planar arrays of semiconductor elements. Planar arrays of semiconductors can be fabricated by lithography at a much lower cost than the cost of fabricating prior art magnetic ferrite memory elements. Additionally, these semiconductor arrays are more compact and faster than prior art ferrite magnetic memory elements. Future benefits of increasingly smaller scale in semiconductor memory are now jeopardized by the concern of loss of reliability, since very small scale semiconductor elements are not electrically robust.\nNon-volatile magnetic memory elements that are read by measuring resistance have been previously demonstrated by Honeywell Corporation. These systems operate on the basis of the classical anisotropic magneto-resistance phenomena, which results in resistance differences when the magnetization is oriented perpendicular versus parallel to the current. Previous work by others has shown that a 2% change in resistance is sufficient to permit the fabrication of memory arrays compatible with existing CMOS computer electronics. Unfortunately, scaling of these elements down from the current 1 .mu.m size has proved challenging."}
{"text": "A Personal Data Vault (PDV) is a cloud, or network, based data storage where a user, or an agent acting on behalf of the user, can store data, such as files, documents, photos, medical journals, music, contact lists, presence information, and so forth. The stored data can subsequently be accessed by authorized clients, i.e., clients which have been granted access to the stored data. The clients may be other users of the data storage, a network operator, or a service provider of, e.g., a personalized service, or a data sharing service.\nIt is known to use policy languages, such as the eXtensible Access Control Markup Language (XACML), to describe how resources, e.g., data or information, are to be treated. Before a resource can be accessed by a client, it is checked if access is allowed by the policy.\nIn situations where the data storage provider, or an associated PEP, is not fully trusted to enforce access policy, encryption is a common solution. However, encryption has the disadvantage that the owner of the data has to be online in order to provide a client requesting access with a decryption key for the encrypted data. If the owner wants to grant offline access, the key has to be provided to the requesting client by some other means than storing the key in the data storage, since storing keys in the data storage defeats the trust model. Regardless how access to the key is implemented, the owner does not have full control and anyone being in possession of the key can access the stored data. Thus, confidentiality of data is not guaranteed anymore.\nA known solution for controlling access to stored data is to employ split keys, e.g., in situations where several clients are available and at least a few of them are trusted. The owner of the stored data may then apply so-called secret sharing of the keys, or possibly of the data, in such a way that only several clients, i.e., more than one client, jointly can retrieve the encrypted data. For example, the owner may split the key among n clients in such a way that at least n/2+1 clients, i.e., a majority of the clients, are required to co-operate in order to obtain the key or the stored data. Several schemes for secret sharing are known in the art\nWhile splitting of keys provides an improved protection, straight-forward usage of this technique implies that a group of clients holding parts of the key can reconstruct the complete key without involving the owner of the data. Also, existing schemes for policy and access control do not provide security if non-malicious mistakes are made. For example, a data storage provider may by mistake grant access to clients not authorized by the owner, e.g., due to misconfiguration of access policies."}
{"text": "1. Field of the Invention\nThis invention relates to the field of electric transformers, and in particular to transformers that are used in ballast circuits of lamp assemblies.\n2. Description of Related Art\nBallasts are commonly used in lamp assemblies to provide a preferred, or optimal, current and voltage to the lamp device. Most lamp ballasts employ a transformer to effect the required transformation of supply voltage to this preferred voltage. The reliability, or expected time-to-failure, of a transformer is inversely proportional to its operating temperature. Electric current flowing through the coils of the transformer generates heat, and this heat causes an increase in the operating temperature of the transformer, thereby reducing its reliability. The amount of heat generated can be reduced by using larger sized wires in each coil, but would result in a larger sized and more costly transformer. Alternatively, the operating temperature can be reduced by efficiently removing the generated heat from the transformer. A variety of techniques are currently available for increasing the efficiency of heat transfer from the coils of a transformer. Thermally conductive potting compounds have been used to conduct the heat from the transformer coils to an enclosure containing the transformer. These semi-fluid compounds, however, are somewhat difficult to handle, compared to solid devices and components, and often introduce reliability problems to other devices, for example, by flowing into the moving parts of switches, relays, connectors, and the like.\nEuropean Patent Application EP 0 254 132, filed Jul. 8, 1987, discloses fastening a metal shell about a transformer, wherein the shell preferably contacts the exterior layer of the coil windings, via a thin insulating layer. This shell is preferably connected to the ballast enclosure, to transfer the heat generated by the coils to the enclosure, both the shell and the enclosure being made of heat-conductive material. This exterior shell must be configured to allow air to circulate within the enclosure, else the thermal efficiency gained by providing the shell will be reduced by the reduction in radiant heat dissipation. Also, the wires that are connected to the coils must be routed through openings in the shell. To assure an efficient thermal contact between the shell and the coil, both the coil dimensions and the shell dimensions must be controlled. A loose fitting shell will be thermally inefficient, and a tight fitting shell may be difficult to fit onto the transformer.\nIt is an object of this invention to provide a method of manufacturing a transformer that provides for efficient heat-transfer from the coils of the transformer. It is a further object of this invention to provide a method of manufacturing a transformer that provides an efficient heat transfer without introducing additional manufacturing constraints. It is a further object of this invention to provide a transformer that includes an efficient heat transfer arrangement.\nThese objects and others are achieved by providing one or more heat-conducting devices coincident with the plastic bobbin that is typically used to form the coils of the transformer. The coil wire is subsequently wrapped around the bobbin and adjacent heat-conducting device assembly, using conventional coil winding techniques. The heat-conducting device is preferably configured to extend beyond the transformer so as to contact a heat-conducting surface, such as a ballast enclosure, when the transformer is appropriately mounted. Because the heat-conducting device is located adjacent the inner windings of the coil, which is typically the locale of the highest temperature build-up in a transformer, a highly efficient heat-transfer is achieved. Preferably, the heat-conducting device is a thick copper conductor having a thin layer of insulating tape separating it from the coil windings."}
{"text": "The invention relates to a surgical cutting instrument for insertion into a cutting trocar sleeve comprising an essentially tubular housing shank, a distal section adapted to be inserted in the trocar sleeve and joined with the housing shank, said distal section having a cutting means and actuating means for the said distal section at a proximal end of the housing shank, the distal section being constituted by a plurality of elements articulatingly joined together, which elements are able to be shifted between a position aligned with the housing shank and an arcuately laterally extended position."}
{"text": "A biaxially stretched polyamide resin film using a polyamide such as nylon 6 and/or nylon 66 is excellent in the mechanical properties such as tensile strength, adhesive strength, pinhole strength and impact-resistant strength, and additionally, in gas-barrier property and heat resistance. Thus, laminated films, in which a biaxially stretched polyamide resin film is used as a front substrate and a sealant made of a polyolefin film is bonded to the front substrate by a method such as dry laminating or extrusion laminating, are used in wide fields including packaging materials for use in sterilization treatment such as boiling or retorting.\nSuch biaxially stretched polyamide resin films are usually used as front substrates, and are in many cases free from direct contact with contents. Accordingly, the behavior of the caprolactam monomer (hereinafter, abbreviated as “monomer” as the case may be) in the biaxially stretched polyamide resin films has not been much mentioned yet.\nIn these years, however, the issue of the deterioration of packaged articles and contents has undergone increasingly severe requests thereto, and the improvement of the issue has come to be demanded. In particular, in the medical applications or the like objecting to subtle compositional changes of the contents, the small molecular weight monomers contained in the polyamide resin film pass through the sealant to migrate into the contents, when heating, for example, for sterilization treatment is conducted, and hence it comes to be impossible to leave such an issue out of consideration.\nFor the purpose of coping with the issue, there have been proposed polyamide resins, in each of which the molecular weight of the constituent monomer unit is large, such as nylon 11 and nylon 12 or copolyamide resins mainly composed of nylon 11 and nylon 12 (JP 4-325159 A). Additionally, a copolymerized polyamide resin between 1,6-hexanediamine and sebacic acid has also been proposed (JP 2001-328681 A). However, these are specific polyamides, and are high in price and low in versatility. Consequently, strongly demanded are films in which highly versatile nylon 6 and/or nylon 66 is used and the amount of the monomer contained therein is low.\nEven if the unreacted monomers and oligomers are removed from a polyamide resin at the stage of being chips prior to film molding, remelting of the polyamide resin chips with a melt extruder or the like regenerates monomers and oligomers, and consequently the monomers remain in the film to degrade the quality of the film. In particular, a polyamide in which caproamide is the main repeating unit thereof has a characteristic that the monomer tends to be more easily generated than in a polyamides formed of a dicarboxylic acid and a diamine.\nIn general, when the terminal group concentration of a polyamide resin is higher, the regeneration amount of the monomer at the time of remelting tends to be larger. Thus, there has been developed a polyamide in which the above-described problem is alleviated by adding a compound capable of reacting with the carboxyl terminals or the amino terminals of the polyamide. Specifically, there has been disclosed a method in which an organic glycidyl ester is reacted with the carboxyl groups and the amino groups of the polyamide (JP 10-219104 A). However, in this method, when the organic glycidyl ester and the polyamide chips are dry blended and melt-kneaded in an extruder, the organic glycidyl ester is allowed to react with the terminal groups of the polyamide. Therefore, in this method, it is difficult to perform uniform mixing in the dry blending step prior to film molding. Consequently, such non-uniform mixing offers a cause for the compositional variation. Thus, it is difficult to obtain a polyamide having a uniform terminal group concentration, and moreover, the dry blending step itself is unsuitable for films involving large melt extrusion amounts. Additionally, the amount of the monomer extracted after the melt molding remains to be as large as 0.35 to 0.5% by mass to show that the reduction amount of the monomer is insufficient.\nOn the other hand, there has been disclosed a method in which the terminal amino groups of a polyamide resin are blocked with a dicarboxylic acid anhydride (JP 2005-187665 A). However, the amount of the regenerated monomer at the time of melting remains to be as large as 0.27 to 0.75% by mass, revealing that it is difficult to sufficiently reduce the amount of the monomer extracted from the polyamide resin film.\nOn the other hand, in these years, recognized is a trend to regulate the discharge, from industrial plants and business institutions, of organic compound materials (generally abbreviated as “VOC”) which evaporate at normal temperature and pressure and easily volatilize into the air. For example, in Japan, on the basis of the revised Air Pollution Control Law, a government ordinance that specifies the type and the size of the institution as an object of regulation came into effect on Jun. 1, 2005. Additionally, the government and ministry ordinances for the discharge standard value, the notification items, the measurement methods and the like were proclaimed on Jun. 10, 2005 and came into effect on Apr. 1, 2006.\nFurther studies are needed as to whether or not the caprolactam monomer discharged in the air provides adverse effects. However, in the production of a polyamide resin film, in printing on the film, and in the steps of laminate processing, bag forming processing and the like using the film, it is the manufacturer's responsibility to reduce the amount of the caprolactam monomer discharged into the air from the film.\nAccordingly, the reduction of the amount of the caprolactam monomer extracted from the film and the recovery of the caprolactam monomer at the time of production of the film are strongly demanded."}
{"text": "1. Field of the Invention\nThe present invention relates to techniques for increasing the efficiency of software on mobiles devices. More specifically, the present invention relates to a method and an apparatus for reducing the number of heap handles in a program, which reduces the effort and resources involved in manipulating the heap handles.\n2. Related Art\nThe Java 2 Platform, Micro Edition (J2ME™), has become a very popular software platform for memory-constrained devices such as wireless devices. Motorola, Nokia, NTT DoCoMo, RIM, Siemens, and many other key players in the wireless device industry have shipped a large number of J2ME™-enabled devices. In fact, over one billion J2ME™-enabled mobile phones have been shipped during the past few years.\nA number of techniques for conserving memory have been developed in order to effectively run applications on such memory-constrained computing devices. One such technique uses exact garbage collection on a shared heap to reclaim memory that is no longer in use. Exact garbage collection uses precise knowledge about all pointers to the heap and between heap objects when scanning and relocating objects. Pointers to the heap are often referred to as “heap roots,” and typically include static pointers and local heap pointers in stack frames.\nExact garbage collection is often implemented using a programming language which is not associated with a garbage-collected heap. For instance, a J2ME™ virtual machine might be written using the C or C++ programming languages. The compilers for such languages are typically unaware of garbage collection, and often do not provide support for locating heap pointers in stack frames. Without such dedicated language support, heap pointers in programs are typically implemented using special structures known as “heap handles.” These heap handles are typically allocated on the stack and require proper management. For instance, heap handles need to be: initialized; inserted into the “root set” of heap roots; and removed after use. Note that the root set is often implemented as a linked list of heap handles.\nUsing heap handles incurs both direct and indirect costs. For instance, additional code needs to be generated and executed to insert and remove heap handles from the root set, and heap handles consume more memory than ordinary pointers in stack frames. Additionally, compilers that are unaware of heap-handle semantics cannot optimize heap handles in the same manner as ordinary pointers, since heap handles are typically at least two-field data structures, and heap handles escape local scope the moment they are added into the root set. Furthermore, if a heap handle is not promptly removed from the root set or set to null after the corresponding object for the heap handle is used, the garbage collector considers the object to be reachable, and does not dispose of the object.\nHence, what is needed is a method and an apparatus that reduces the above-described costs of manipulating heap handles where possible."}
{"text": "1. Field of Invention\nA trailer mounted shooting bench, storage container and target apparatus provides a portable shooting station, which can be transported to a remote location appropriate for target shooting or hunting, and provides the target attaching to the trailer mounted shooting bench during transport with the target being detached and deployed in a desired location with the trailer mounted shooting bench moved to a separate location, the shooting bench made stationary for targeting and shooting at the target, or used for any general shooting purpose. The shooting bench provides various storage for shooting paraphernalia with an array of options provided for sport shooting or hunting purposes.\n2. Description of Prior Art\nA preliminary review of prior art patents was conducted by the applicant which reveal prior art patents in a similar field or having similar use. However, the prior art inventions do not disclose the same or similar elements as the present portable shooting bench and target, nor do they present the material components in a manner contemplated or anticipated in the prior art.\nPortable shooting benches found in prior art reveal several portable shooting devices which provide a stable surface from which to shoot. In U.S. Pat. No. 5,284,280 to Stonebraker, Sr., a backpack shooting table is disclosed which provides a four legged backpack frame component with independent adjustable legs, a stable horizontal platform upon the frame with an adjustable gun rest. A collapsible and portable gun table with four adjustable and collapsible legs is demonstrated in U.S. Pat. No. 5,697,180 to Morizio, which is approximately the size of a card table. This table also provides an adjustable gun rest, drawers, transport wheels and can be used as a table or an angled upright support. A portable shooting bench assembly is further disclosed in U.S. Pat. No. 7,549,247 to Reese, which provides an adjustable gun rest and a bench seat for seated shooting, the assembly being presented with a wheelbarrow-type frame with two handles near the set for picking up the seat end and a large single wheel near the gun rest end for rolling the assembly. It too has adjustable legs and further provides for seat adjustment for a left handed or right handed shooter.\nSeveral other shooting related devices are presented in prior art that incorporate into a rear vehicle trailer hitch. In U.S. Pat. No. 6,684,550 to Highfill, a clay target thrower is mounted to an arm that extends up and away from a receiver hitch. A seat attached to a horizontal arm having a single adjustable jack and an end defining an adjustable height table top, as shown in U.S. Pat. No. 6,269,578 to Callegari, is connected into a receiver hitch. A near identical receiver hitch mounted device is disclosed in U.S. Pat. No. 6,935,064 to Thompson, except it has a seat that swivels from one side to another and a table top which may be presented for a left handed shooter or a right handed shooter. A free standing shooting table with several horizontally presented surfaces for elbow rests, gun stock rests and other swivel mounted items, including a lower seat, having independent collapsible legs for support, with this device being transported in a receiver hitch when collapsed and partial disassembled. In U.S. Pat. No. 7,536,820 to Wade, a receiver hitch mounted shooting rest is shown with multiple adjustments for height and distance from the tailgate are provided, while U.S. Pat. No. D628,405 to Pippin merely shows a design patent for a receiver hitch mounted shooting table with a vertical height adjustment and foot rest pegs."}
{"text": "The availability of new technologies has given vending machine manufacturers and software developers many tools to address market demands of vending operators. Advances in software and electronics are now enabling the use of computer controls and data acquisition systems directly inside a vending machine. Some of the latest vending machines now make it possible for vending machine operators to download selected aspects of operational information on-site onto portable computers. Although these computerized systems make it easier for operators to gather and analyze some data, they generally ignore many significant aspects of vending operations and commonly provide untimely data. In addition, these computerized systems are typically cumbersome, difficult to connect and fail to leverage advanced functionality likely to enhance vending machine operation efficiency and profitability."}
{"text": "The existing state of the art discloses different types of methods for modeling an object. Said methods are mainly classified into passive methods and active methods. In the area of active methods, sensors such as laser or structured light scanners or also Time-of-Flight type cameras are used. There are other possibilities such as projecting, with the aid of a video projector, a known pattern on an object and deducing the shape of the object by means of analyzing the deformation experienced by the pattern due to the shape of the object.\nIn the area of passive methods, most techniques exploit the geometric triangulation relating two or more views of the object of interest.\nThe present invention is in the field of the passive methods which include the following approaches of modeling objects based on views:                Structure from Motion, SfM, consisting of estimating the model of the scene in front of a camera in motion. However, the technique is only applicable to a set of multiple static cameras. Generally, an SfM algorithm establishes the match between the views of a set of points in the scene. By means of establishing this match, it is possible to triangulate the position of the points in the three dimensions of the space in front of one or several cameras. From this point, there are several possibilities for generating a model of an object. One possibility is using triangulation to calibrate the position of the camera throughout its motion or the position of each static camera. A dense model of the shape of the scene can be obtained, for example, by means of Shape from Stereo. Another possibility is assuming that the surface between any three points is locally flat. This model is therefore obtained by connecting points in groups of three by a triangle. The set of 3D triangles form a mesh representing the shape of the object. In this sense methods which reconstruct parts of the flat object are known from the state of the art. Firstly, matches between flat segments are established. Four points per segment or region are found and then a homography is induced. This homography allows establishing the epipolar geometry between the views. Finally, the set of segments can be positioned in 3D.        3D volumetric reconstruction. This approach encompasses from the least to the most precise modeling. For example, the box delimiting the real object would be a too coarse model. There are more precise models such as the Convex Hull (CH), the Visual Hull (VH) and the Photo Hull (PH). One of the most widespread volumetric models due to its good ratio between precision and low computational cost is the Visual Hull (VH). The Visual Hull is obtained by means of a method referred to as Shape-from-Silhouette (SfS). In a first phase, the Shape-from-Silhouette extracts the active entities of the scene (silhouettes of the object) by means of a set of cameras. The Visual Hull therefore corresponds with the volume inside the intersection of the cones going from the optical center of the cameras through the silhouettes in the optical planes of the cameras. The set of cameras must be intrinsically and extrinsically calibrated beforehand. The calibration can thus be obtained using the set of control points the coordinates of which are automatically known as a set of characteristic key points, as in the Structure from Motion approach.        Shape from Shading, “SfSh”, deals with recovering the shape from a gradual variation of the shading in the view. The idea behind Shape from Shading is that the color intensity can be described as a function of the surface, the shape, and the direction of the light source. Most SfSh algorithms assume that the direction of the light source is known.        \nThe passive methods described above have several drawbacks depending on the method used. In the case of methods based on Structure from Motion (SfM), the drawbacks arise from the objects without texture. In fact, in the absence of texture on the surface of the object, the resulting model is very coarse. In the case of very limited texture but with sufficient points for calibrating the set of cameras, the Shape from Stereo method can be used. However, the result of the previous method has the drawback that it is not capable of isolating the object from the objects forming the background or surrounding the object which is being modeled. In the particular case of the methods described above and which are based on finding four points of a segment and generating a homography, the entire calibration process depends on the possibility of establishing a match between the detected planes, which is not viable for objects without texture.\nOn the other hand, the Visual Hull obtained with a generic SfS method mainly depends on two aspects. Firstly, the positions of the cameras determined the efficiency of the SfS method. Another limitation of the applicability of this method is that the silhouettes are extracted by comparison with a known static background. This means that the object cannot be present in the scene when the background is captured. Consequently, this method is only valid for objects which can be easily obtained or introduced in the scene but not for modeling a part of a room, such as a wall or a fixed board.\nIt would therefore be desirable to find a method for generating a model of a flat object from views of the object which does not depend on the texture of the object to be modeled or on the consequent limitation involved in the correct calibration of the cameras, as well as to the capacity to move the object to be modeled with respect to the background or to the site in which the object is located."}
{"text": "In known systems electrical transmission lines leading from a power source such as a generator, to consumers are protected against insulation failure or overload by at least one circuit breaker. In certain instances the circuit breaker includes mechanical switching devices having a pair of conductor terminals and a bridging member for opening and closing the gap in between said terminals. Because it is not possible to interrupt a high voltage or a large electrical current instantaneously, the electric arc emerging in the expanding gap upon pulling the conductor terminals apart is often spread and broken in an insulation gas environment, such as pressurized air or sulfur hexafluoride for example. The high voltage circuit breaker market is increasingly dominated by self-blast technology.\nFR 2575594 discloses a representative of such a self-blast-type circuit breaker (GCB) using SF6 as an extinguishing agent. The arrangement includes movable and immovable electrical contacts located in an arcing zone such that an electric arc is generated in the arcing zone. A pressure chamber arrangement is connected by channels to the SF6-filled arcing zone for enhancing the breaking quality by preventing the electric arc from becoming revitalized after an initial extinction.\nIn known systems, the highest short-line fault ratings (SLF) are covered by puffer type gas circuit breakers such as tank SF6 puffer circuit breakers for example. If limits above 50 kA, 245/300 kV are to be achieved by employing such puffer type circuit breakers expensive line to ground or grading capacitances are specified.\nThere have also been attempts in scaling-up known self-blast technology puffer breakers to withstand ratings of 63 kA at 300 kV, in a 60 Hz environment having 450 Ohm without a delay in time.\nKnown GCB features a quenching chamber, also known as interruption chamber, which is filled with an insulating gas. The chamber extends along a longitudinal axis and is designed to be radially symmetric, (e.g., rotationally symmetric about said longitudinal axis. The quenching chamber further includes at least two separable arcing contact pieces coaxially arranged and facing each other as well as an arcing zone formed in between the at least two arcing contact pieces. An electric arc burns between the at least two arcing contact pieces during a disconnection/interruption process and heats up the isolating gas in the arcing zone when the contact pieces are separated. The heat causes an increase of the pressure of the insulating gas in the arcing zone of the GCB. The pressurized gas escapes through at least one dedicated annular gap between an arcing contact piece and the quenching chamber and through cavities arranged proximal to the longitudinal axis in the contact pieces, if any, such that each emerging flow path constitutes an optimal gas nozzle. Thus, in the context of the present disclosure, the term nozzle describes a functional flow rather than a physical component.\nKnown attempts to achieve the above ratings with scaled-up self-blast technology puffer breakers failed because higher pressure values are expected which lead to mechanical failure of the material of the GCB and an undesired reduction of the dielectric withstand of the insulating gas due to the associated high temperature of above 2000K.\nThere are two situations under which a high-voltage circuit breaker, in particular a high-voltage alternating current circuit breaker, should endure. The first situation is known as a short line fault (SLF) and the second situation is known as a terminal fault (T100a).\nIn a GCB, the pressure in the arcing zone should be comparatively high for extinguishing the electric arc in a reliable manner in case of a short line fault. One drawback, however, is that a high pressure raises the thermal load to the structure of a circuit breaker. With respect to a terminal fault, the current pressure values in the arcing zone exceed the pressure values that are specified for reliably extinguishing the electric arc, which are comparatively low. Hence, in case of a GCB, the gas nozzle should be able to bear the pressure in the arcing zone in the SLF situation, as well as withstand T100a conditions.\nIn “Investigation of Technology for Developing Large Capacity and Compact Size GCB” disclosed in the IEEE Transactions on Power Delivery, Vol. 12, No. 2 dated April 1997, a different solution for achieving the above-mentioned ratings by employing different nozzle geometry is proposed. This nozzle is different from other known GCB applications because of an inner nozzle that is assigned to a movable arcing contact wherein said inner nozzle contributes to the establishment of local higher gas pressures specified for the thermal interruption at a SLF without only increasing the pressure in a dedicated puffer chamber of a GCB.\nThere remains the drawback, that high gas pressures are known to cause high temperatures which in turn are undesired for dielectric interruption since the gas becomes conductive above 2000 Kelvin such that it can not be employed sensibly for breaking an electric arc in case of SF6 gas employed as the extinguishing agent in a GCB."}
{"text": "1. Field of the Invention\nThe present invention generally relates to brassiere. More particularly, the present invention relates to brassiere, especially sport brassiere, having breast supporting panels in the breast cups.\n2. Description of the Prior Art\nA brassiere attempts to combine comfort for the wearer and support for the breasts, which can be mutually exclusive features. Commonly, a brassiere incorporates stretchable or elastic materials for the comfort of the wearer. Typically, support for the breasts is achieved with underwires made of rigid materials, such as metal and/or plastic.\nThe difficulty in making both a comfortable and supportive brassiere is amplified for xe2x80x9csportsxe2x80x9dbrassiere designed to be worn during exercise and other physical activities. Some sports brassiere function by compressing the breasts to the body of the wearer. A disadvantage of this type of sports brassiere is that it does not support the breasts of the wearer and the breasts may shift within the sports brassiere during exercise and other physical activities. To reduce the risk of injury to the breasts and back, sports brassiere are expected to be strong and fully supportive. At the same time, sports brassiere are sometimes worn for extended periods of time. Therefore, such brassiere must be quite comfortable to wear. Also, for various reasons, some women prefer to wear sports brassiere exclusively, instead of everyday and/or xe2x80x9cfashionxe2x80x9d brassiere.\nIn light of the foregoing, there is an ongoing need for brassiere, especially sports brassiere, that are both comfortable to wear and adequately support the breasts of the wearer.\nIt is an object of the present invention to provide a brassiere that is comfortable to wear.\nIt is also an object of the present invention to provide a comfortable brassiere that can adequately support the breasts of a wearer during exercise and other physical activities.\nIt is a further object of the present invention to provide such a comfortable and supportive brassiere that is simple to manufacture.\nIn light of the foregoing, there is provided a brassiere having a body made of a stretchable material for comfort. Each breast cup has a breast receiving portion with an inner layer and outer layer, and a support panel portion with the inner and outer layers and a breast supporting panel between the inner and outer layers. The breast supporting panel is substantially inelastic in all directions and, preferably, forms an arcuate edge for supporting the breast. Thus, the entire brassiere stretches and moves with the wearer, while each non-elastic breast supporting panel provides the desired breast support. Significantly, the bottom edge of each breast supporting panel is spaced from the bottom edge of the brassiere. Decreasing the length of the breast supporting panel so that its bottom edge does not contact the bottom edge of the brassiere allows the bottom periphery of the brassiere to have additional flex that improves the overall comfort of the brassiere. Preferably, the bottom edge of the brassiere is turned to form the waistband of the brassiere."}
{"text": "In prior U.S. Pat. No. 5,052,627 of the present inventor issued Oct. 1, 1991 is disclosed a machine for spreading particulate material and particularly, but not exclusively, granular fertilizer in an agricultural situation. This machine includes a tank arranged to be mounted on a suitable transportation vehicle. At one end of the tank which is either the front end or the back end is mounted a pair of booms with each boom extending outwardly to a respective side of the tank for movement with the vehicle across the ground. Each boom includes a series of pipes which are horizontal and arranged either in a single row or in two rows so that the pipes are connected side by side. The tank has a pair of belts extending longitudinally of the tank side by side in a common horizontal plane. Each belt feeds out of the end of the tank for discharging material carried on the upper run of the belt into a guide duct system positioned underneath the end of the belt. The guide duct system includes a plurality of ducts equal to the number of pipes and arranged side by side across the width of the discharge end of the belt. The guide duct system thus receives portions of the material from different positions across the width of the belt and transfers that portion of material to a respective one of the pipes. The inner end of the pipes are connected to a manifold so that air from a fan blows through the pipes to carry the portion of material along the respective pipe to a spreader system at the end of the pipe. The spreader system can in some cases include two separate spreaders one directly at the end and one spaced inwardly from the end so as to split the material at the two separate spreaders thus reducing the number of individual pipes.\nThis machine has achieved considerable commercial success and is widely sold in North America.\nThe above patent also discloses a system for adding one or more additional components to the first component from the tank. Thus the belt carries the first component out of the tank using a gate at the exit of the tank for levelling the material on the belt at a predetermined thickness. Two further tanks are provided each of which has an individual metering system which meters an additional component from the additional tank onto the top of the same belt so that the same belt carries the additional material with the first material from the main tank to the guide system for discharge of the portions into the separate ducts of the guide system.\nHowever the machine described above using the supplemental tanks does not provide sufficient flexibility of the materials to be supplied by the belt to accommodate the variations in the components of the fertilizer which are desired for modern agriculture."}
{"text": "1. Field\nEmbodiments relate to a secondary battery and a secondary battery module.\n2. Description of the Related Art\nSecondary batteries are rechargeable batteries, and may be broadly used in portable electronic devices, e.g., cellular phones, notebook computers, and camcorders.\nA secondary battery is formed by inserting an electrode assembly, in which a positive electrode, a negative electrode, and a separator are wound in the form of a jelly roll, into a case through an opening of the case, and covering the opening by using a cap assembly. Current collecting plates are formed at two ends of the electrode assembly and are electrically connected to a terminal unit of the cap assembly. Accordingly, if an external terminal is connected to the terminal unit, current generated in the electrode assembly may be provided to the external terminal via the current collecting plates and the terminal unit.\nThe terminal unit may include a positive electrode rivet and a negative electrode rivet connected to the current collecting plates, and rivet terminals bonded to the positive and negative electrode rivets in order to be connected to bus bars. Bonding between the positive electrode rivet or the negative electrode rivet and a rivet terminal, and between a rivet terminal and a bus bar may be performed by using a laser welding method. However, since the positive and negative electrode rivets are generally formed of dissimilar metals, if the rivet terminals are formed by using one metal, dissimilar metal welding may be performed between the positive electrode rivet or the negative electrode rivet and a rivet terminal."}
{"text": "A comminuting apparatus may typically be used for example for comminuting wood, paper, plastic material, rubber, textiles, production residues or waste from trade and industry, but also for dealing with bulky refuse, domestic refuse, collections of paper and other waste materials for example from organisations set up to dispose of waste and such like in an environmentally friendly fashion, as well as more specialist waste such as hospital and clinical waste. A comminuting apparatus for such a purpose may comprise a drive unit with at least one electric motor having a drive shaft operatively connected to a comminuting shaft. At its periphery the comminuting shaft has comminuting tools over its working width. The tools co-operate with a counterpart means adapted in respect of shape to the rotational surface of the comminuting shaft, for comminuting the material to be processed. In such an apparatus the material to be comminuted is comminuted by cutting, shearing, squeezing, tearing and/or rubbing, between rotor members or in the co-operation between a rotor member and a fixed transverse member operatively associated therewith. Such an apparatus may be found for example in EP 0 419 919 B1.\nThere are also forms of comminuting apparatus comprising a plurality of rotors each with a respective stationary transverse member associated therewith, between the respective rotors.\nTo perform an operation of roughly pre-comminuting waste material, a rotary speed of the comminuting shaft of between about 20 and 50 rpm is appropriate. Hydraulic drives are generally used for that purpose. When dealing with material which is easy to comminute or which is already sufficiently pre-comminuted, such as for example films, sheets, packaging residues and the like, the comminuting apparatus can in principle be operated at higher rotary speeds in order to increase the waste material throughput, and in that respect presentday comminuting apparatuses are equipped with comminuting shafts which may be driven at between about 80 and 500 rpm. The electrical drive power of such an apparatus is between about 30 and 450 kW.\nVarious drive configurations may be adopted for such comminuting apparatuses. Conventional apparatuses generally include an asynchronous motor which is preferably of a 4-pole configuration and which accordingly operates at a motor speed of 1500 revolutions at a mains frequency of 50 Hz. To set the specified speed of rotation of the comminuting shaft, the transmission of force thereto from the motor is effected by way of a belt drive or a universally jointed shaft or a clutch to a transmission in which the rotary speed, depending on the respective demands involved, is reduced to between about 90 and 200 rpm, whereby the torque at the comminuting shaft is increased in comparison with that of the motor in the same relationship.\nIn regard to a further design configuration of a comminuting apparatus, it has a drive in the form of an electric motor which is generally of a 4-pole or 6-pole design and which accordingly operates at 1500 rpm or 1000 rpm respectively at a mains frequency of 50 Hz. Connected downstream of the electric motor is a transmission operating with a pulling means such as a belt or chain transmission. That arrangement makes it possible to attain rotary speeds for the comminuting shaft of between about 200 rpm and 500 rpm, by means of a simple drive, although it will be noted that belt pulleys which are very large and usually expensive have to be used.\nAs the large belt pulleys employed have a high moment of inertia, a load-limiting or load-separating clutch or coupling unit is generally fitted at or in a hub between the comminuting shaft and the belt pulley, preferably a slipping clutch, depending on the material to be comminuted, in order to avoid breakage of the comminuting shaft. At even lower rotary speeds, it is necessary to use a double-run belt transmission. In that case, very high levels of torque can be produced at the comminuting shaft, which however require suitable dimensioning of the drive elements, so that such a design configuration is very expensive and maintenance-intensive, while at the same time the comminuting apparatus takes up a great deal of space, by virtue of its bulky structure. The fluid couplings which are generally used in both the above-discussed drive configurations optimise the known disadvantageous start-up characteristic of an asynchronous motor and make it easier for the comminuting shaft to start under load. In addition, in the event of a sudden blockage, for example due to the presence of a foreign body in the material being comminuted, the coupling arrangement has a damping effect and reduces the load peaks which are produced by the apparatus in the supply mains network.\nA further conventional drive arrangement for a comminuting apparatus employs an asynchronous electric motor, a hydraulic pump and an oil motor. The moment produced by that drive assembly is passed to the comminuting shaft with or without an interposed transmission. That design configuration is highly expensive and maintenance-intensive, and comparatively unfavourable in terms of level of efficiency, while in addition the apparatus is very noisy. On the other hand that configuration affords the advantage that the rotary speed of the comminuting shaft can be adjusted over a predetermined range.\nWhat is common to all those conventional drive configurations is that they include a plurality of drive members for connecting a motor to a comminuting shaft. They are comparatively expensive, they increase the amount of space required and in addition increase the level of noise generated by the apparatus. Connecting a plurality of drive members in succession results in the machine suffering from a power loss. In other words, the machine has an unfavourable level of efficiency, with a corresponding energy loss. As the entire drive consists of a plurality of drive members, those drive members in combination exhibit a high level of mass moment of inertia, which, in the event of load peaks which suddenly occur, can result in problems in regard to strength and operating life and under some circumstances can result in parts of the machine being broken and destroyed. Load peaks of that kind can occur on the one hand due to pieces of material which cannot be comminuted, for example metal, stones, rocks and so forth, in the material being processed, but they can also occur when comminuting tough resilient materials with a high level of tearing strength such as for example fiber mesh or web, cables, cords and the like.\nDepending on the material being comminuted, the rotor blades adopted and the rotary speed of the rotor or rotors, rotary oscillations often occur, in particular when gear transmissions are used in the drive assembly. Such oscillations generate a large amount of noise and reduce the service life of the drives."}
{"text": "Semiconductor manufacturers have developed components, such as packages and BGA devices, which contain multiple semiconductor dice. For example, systems in a package (SIP) include multiple dice having different configurations, such as a memory, a processing, or an application specific configuration. The multiple dice provide increased integration, security and performance in a component.\nOne aspect of these multi-dice components is that they typically have a relatively large peripheral outline and thickness. For example, conventional systems in a package have two or more dice spread out on a common substrate. These components are typically larger than conventional plastic semiconductor packages. It would be desirable to be able to fabricate semiconductor components, such as packages and BGA devices, with multiple dice, but also with a chip scale outline and thickness.\nAt the same time, components need a reliable and efficient internal signal transmission system, and a high input/output capability. One aspect of conventional chip scale components, such as chip scale packages (CSP), is that they are difficult to manufacture with the reliability required in the industry. For example, some chip scale components include relatively complicated signal transmission systems, such as beam leads and wire conductors. These signal transmission systems are difficult to manufacture, and are prone to failure, particularly at the high pin counts required for demanding electronics applications. It would be desirable for a multi-dice component to have a reliable signal transmission system capable of volume manufacture.\nThe present invention is directed to a multi-dice component having a chip scale outline, an integrated internal signal transmission system, and a high input-output capability. In addition, the present invention is directed to wafer level methods for fabricating multi-dice, chip scale components."}
{"text": "Radiofrequency (RF) powered electron beam accelerators (or accelerator guides) have found wide usage in medical accelerators where the high energy electron beam is employed either directly for therapeutic purposes, or converted to generate x-rays for therapeutic and diagnostic purposes. The electron beam generated by an electron beam accelerator can also be used directly or indirectly to kill infectious pests, to sterilize objects, and to change physical properties of objects and materials. A further common use of electron beam accelerators is to perform radiographic testing and inspection of objects, such as containers for storing radioactive material, and concrete and steel structures.\nThe RF power for an electron beam accelerator is generally desired to be controlled, such that the beam energy from the accelerator can be delivered in a desired manner. It is common practice that the RF power be delivered to the accelerator as a series of short pulses, resulting in an electron beam output of a corresponding series of beam pulses. In some applications, it may be desirable that the accelerator be capable of generating beam energy pulses that vary between different energy levels, even on a pulse-by-pulse basis. However, existing systems may not be able to accomplish these objectives. Also, existing RF systems may not be able to control generated power such that power delivered to the accelerator can be varied quickly, e.g., in the order of milliseconds, between at least two power levels, which may be desirable in certain accelerator system applications.\nFurther, in existing systems, RF power provided by a power generator to an accelerator may be reflected back to the power generator. In many applications, it is desirable that such reflected RF power from the accelerator be controlled such that the frequency of a power generator will be “pulled” to the accelerator frequency, resulting in a stable operation of the power generator and the accelerator. If the reflected power is not controlled, the frequency of the power generator may be forced or “pulled” away from the operational frequency of the accelerator, resulting in failure of the accelerator to operate correctly."}
{"text": "1. Field of the Invention\nThe present invention relates to a horizontal axis wind turbine and a method for measuring an upflow angle.\n2. Description of Related Art\nIn recent years, horizontal axis wind turbines have been in practical use in order to gain electric power from natural wind. The performance of the horizontal axis wind turbines is shown by a power curve representing relationship between the wind speed at hub (the wind speed at the rotational central portion of a rotor of a horizontal axis wind turbine) and power (production of electricity) in general. The production of electricity leading directly to profitability is predicted based on the power curve and the wind speed at hub estimated by simulation or observed in advance.\nWind-generated electricity systems are often installed in complex topography. In such topography, the wind speed varies with not only the height from ground but also the horizontal position, and in addition, an upflow wind is often generated. The power generation and structure damage are affected by not only horizontal component of the wind but also vertical. Therefore it is meaningful to measure 3-D wind speed with consideration of “upflow angle” for improvement of predictability of the production of electricity.\n3-D ultrasonic anemometers and Pitot tubes are proposed as earlier developments for measuring 3-D wind speed with consideration of a “upflow angle”. (For example, see Kaijo Corp. “measurement and control system business—atmospheric apparatus” [online] 1997, Kaijo Corp.)\nThe above-described 3-D ultrasonic anemometers, however, have a problem that it is expensive and large. An anemometer for measuring the wind speed at hub of a horizontal axis wind turbine requires durability to endure even in a relative harsh environment for long periods without maintenance, whereas 3-D ultrasonic anemometers and Pitot tubes are not produced on an assumption of such an operating environment. Therefore they have reservations about durability and lack reliability."}
{"text": "Machine-to-Machine (M2M) communication (also referred to as “machine-type communications” or “MTC”) may be seen as a form of data communication between entities that do not necessarily need human interaction.\nM2M communication may be used in a variety of areas such as: security, tracking/tracing, healthcare, remote maintenance/control and metering. M2M communication may be used in surveillance systems, order management, gaming machines and remote monitoring of vital signs. M2M communication may be used in programmable logic controllers (PLCs), sensors, lighting, vending machine control and in applications related to power, gas, water, heating, grid control, and industrial metering. Additionally, M2M communication based on machine type communication (MTC) technology may be used in areas such as customer service.\nDepending on its implementation, M2M communication may be different from some current communication models. For example, M2M communication may a large number of WTRUs, and/or may involve very little traffic per WTRU. Additionally, relative to some current technologies, M2M communication may involve lower costs and less effort to deploy.\nM2M communications may take advantage of deployed wireless networks based on Third Generation Partnership Project (3GPP) technologies such as Global System for Mobile Communications (GSM), Universal Mobile Telecommunications System (UMTS), Long Term Evolution (LTE), and/or other technologies such as those developed by the Institute of Electrical and Electronics Engineers (IEEE) and 3GPP2. M2M communications may use networks based on these technologies to deliver business solutions in a cost-effective manner. In a circumstance involving ubiquitous deployment of wireless networks, the availability of the wireless networks may facilitate and/or encourage the deployment and use of M2M WTRUs. Additionally, further enhancements to these technologies may provide additional opportunities for the deployment of M2M-based solutions."}
{"text": "A flexographic plate is a type of relief printing plates and generally includes an elastic relief plate made of rubber or photosensitive resin, to which a liquid ink is applied for printing. Flexographic plates can print on rough or curved surfaces and are thus widely used for printing images on wrapping, magazines, cardboards, labels,  and bottles. Flexographic plates were previously manufactured by pouring molten rubber in a mold and then curing the rubber or by manually carving a rubber plate. Neither of these techniques was suitable for producing accurate flexographic plates, however. Lately, the development of a new technique that uses curable resins to make flexographic plate materials has made the production of flexographic plates considerably simple.\nA typical flexographic plate material of the newly developed type includes, from top to bottom, a surface protective layer; a layer of a curable resin composition that is curable by irradiating with an active energy ray and is composed of an elastomer, such as urethane rubber, butyl rubber, silicon rubber and ethylene propylene rubber, an ethylenic unsaturated compound and, if necessary, a photopolymerization initiator; an adhesive layer; and a substrate (See, for example, “Kankosei jushi no kiso to jitsuyo (Basics and applications of photosensitive resins)” Supervised by Kiyoshi Akamatsu, CMC Co. Ltd. (2001) 152-160).\nIn one process for producing a flexographic plate from such a flexographic plate material, a film carrying a negative image of a letter, diagram, picture, pattern, or any other image to be printed is first applied to the surface of the protective film opposite to the substrate. The negative film is then irradiated with an active energy  ray from above, so that the predetermined areas of the curable resin composition layer are selectively cured by the action of the active energy ray transmitted through the imaged area of the film and become insoluble to solvent. Subsequently, the negative film and the protective film are removed and a solvent is applied to remove the non-irradiated or uncured areas of the curable resin composition layer (development step) and thereby form an image area (i.e., image plate surface). This completes a flexographic plate (See, for example, “Kankosei jushi no kiso to jitsuyo (Basics and applications of photosensitive resins)” Supervised by Kiyoshi Akamatsu, CMC Co. Ltd. (2001) 152-160; Japanese Patent Publication No. S55-34415; U.S. Pat. No. 4,323,636; Japanese Patent Publication No. S51-43374; and Japanese Patent Application Laid-Open No. H2-108632).\nIn an effort to ensure formation of fine dots and lines on the flexographic plates and prevent chipping of the image plate surface during development, improvements have been made as to the type and proportion of the resin to be added to the curable resin composition. One example involves the use of a styrene-based block copolymer in which the part of the copolymer formed of a conjugated diene has a significant bound vinyl content (See, Japanese Patent Application Laid-Open No. H5-134410). In another example, a certain thermoplastic elastomer composed of a  monovinyl-substituted aromatic hydrocarbon and a conjugated diene is used in conjunction with a diene-based liquid rubber that has a high average proportion of bound vinyl units (See, Japanese Patent Application Laid-Open No. 2000-155418).\nThe technique described in Japanese Patent Publication No. S55-34415 employs crystalline 1,2-polybutadiene in conjunction with a polymer compound, such as polyisoprene rubber, that comprises as its constituents at least one of ethylene, butadiene and isoprene. A drawback of this technique is that the uncrosslinked rubber used in the resin makes the flexographic plate susceptible to deformation (cold flow) during storage or transportation of uncured plates. U.S. Pat. No. 4,323,636 and Japanese Patent Publication No. S51-43374 describes the use of a certain block copolymer (preferably a styrene-isoprene-styrene triblock copolymer or a styrene-butadiene-styrene triblock copolymer having a particular composition) that is a thermoplastic elastomer and in which the hard segments have a grass transition temperature of 25° C. or above. In these techniques, the part of the copolymer formed of polystyrene causes a cohesive force, which reduces deformation of uncured plates. However, the polystyrene blocks of the elastomer do not undergo crosslinking even when irradiated with an active energy ray, so that the uncrosslinked polystyrene  blocks causes a poor solvent resistance of the cured area. Consequently, the image area tends to swell when a solvent is applied to remove the uncured area, resulting in insufficient reproducibility and poor ink resistance of the flexographic plates. Another flaw of this technique is that the strength and extension of the cured area are insufficient especially for forming a fine image area and the resulting flexographic plate becomes less durable. This causes chipping in the edge of the image plate surface during the removal of the uncured area by washing with a solvent and, when necessary, a brush. As a result, the desired sharp image plate surface may not be obtained.\nIn the technique described in Japanese Patent Laid-Open Publication No. H2-108632, the flexibility of flexographic plates is increased by the use of a binder (preferably, a styrene-butadiene-styrene triblock copolymer) containing thermoplastic and elastomeric domains, in combination with a particular addition polymerizable ethylenic unsaturated monomer. Despite its improved flexibility, the resin according to this technique includes some part formed of a polystyrene block similar to the one described above, which makes the flexographic plate less resistant to solvent. The cured area of the flexographic plate obtained by this technique is not strong enough to form a fine image area.\nAlthough both of the techniques described in  Japanese Patent Application Laid-Open No. H5-134410 and No. 2000-155418 have managed to improve the curability of the part of the styrene-based thermoplastic elastomer formed of conjugated diene units and have managed to increase the toughness of the resulting flexographic plate, the similar polystyrene block part present in the thermoplastic elastomers suppresses the solvent resistance and the flexographic plates are not operative enough to form a fine image area.\nAccordingly, it is an object of the present invention to provide a curable resin composition suitable for the production of a flexographic plate material that allows printing on an article with rough surfaces, such as cardboard and recycled paper. The curable composition of the present invention can be cured to form strong and extendable areas and can thus be used to make flexographic plates that can form a sharp image plate surface even for a fine image. It is also an objective of the present invention to provide a flexographic plate material that uses the curable resin composition as its constituent."}
{"text": "Wireless devices typically are able to send and receive information about the caller placing a phone call. Wireless devices can include, for example, cellular devices, tablets, smart phones, and the like. Stored in the memory, wireless devices often maintain a contact database. The contact database can include a telephone number correlated to a name, address, and company, for example. These contact databases can associate phone numbers with callers. When a recognized caller whose phone number is stored on the wireless device of a called party calls the called party, the called party's wireless device may communicate the caller's information on the called party's wireless device so that the called party can determine whether or not to answer the call. Sometimes, however, the phone number of the caller may not be associated with any contacts stored on the called party's wireless device, so a called party may not be able to determine who is calling. This results in some important calls not being answered or some unimportant calls being answered. On the other hand, answering such anonymous phone calls can result in being subject to undesirable solicitation phone calls and the like.\nTherefore, there is a need for providing a calling party's information when that information is not stored on a called party's wireless device so that the called party's wireless device can communicate that information for the called party."}
{"text": "Computer graphics display systems typically comprise a frame buffer memory which stores the color and Z coordinate associated with each pixel to be displayed on the monitor of the computer graphics display system. A frame buffer controller of the computer graphics display system controls the process of writing the Z coordinates and the colors of the pixels to the frame buffer memory. In many high-performance computer graphics display systems, Z buffer depth comparison tests are used to determine whether a new Z coordinate received in the frame buffer controller corresponds to a pixel that will be visible when displayed, or whether the pixel associated with the new Z coordinate will be occluded or hidden if displayed. If the pixel will be occluded, it is unnecessary to write the Z coordinate and the associated color to the frame buffer memory and the pixel can be discarded. On the other hand, if the pixel will be visible, the Z coordinate and the associated color must be written to the frame buffer memory.\nIn these types of systems, Z buffer depth comparison tests are performed by reading the old Z coordinate for the pixel from the Z buffer, comparing the old Z coordinate with the new Z coordinate, and, if the new Z coordinate passes the Z buffer depth comparison test, writing the new Z coordinate and the associated pixel color into the Z buffer and image buffer, respectively, of the frame buffer memory. Once these steps have been performed, the next pixel is processed in an identical manner.\nAlthough this type of depth comparison test is needed to occlude hidden surfaces, performing depth comparison operations for each new Z coordinate received in the frame buffer controller requires significant memory bandwidth for reading the Z coordinates from the Z buffer.\nAccordingly, a need exists for a method and apparatus for performing Z buffer depth comparison tests which reduce the number of reads to the frame buffer memory and thereby improve memory bandwidth efficiency."}
{"text": "A great effort has been made to develop technologies for cast molding of hydrogel contact lenses with high precision, fidelity and reproducibility and at low cost. One of such manufacturing technologies is the so-called Lightstream Technology™ (CIBA Vision) involving a lens-forming composition being substantially free of monomers and comprising a substantially purified prepolymer with ethylenically-unsaturated groups, reusable molds, and curing under a spatial limitation of actinic radiation (e.g., UV), as described in U.S. Pat. Nos. 5,508,317, 5,583,463, 5,789,464, and 5,849,810. The Lightstream Technology™ for making contact lenses have several advantages. First, the curing process is fast, at a time scale of seconds. Fast curing can ensure design and adaptation of a high speed, continuous and automatic lens production involving on-line lens curing. Second, by using a composition comprising a prepolymer and being substantially free of monomers, subsequent extraction steps (removing unpolymerized monomers from the lenses) required in a traditional cast-molding manufacturing process are eliminated. Without lens extraction, the production cost can be reduced and the production efficiency can be further enhanced. Third, reusable quartz/glass molds or reusable plastic molds, not disposable plastic molds, can be used, because, following the production of a lens, these molds can be cleaned rapidly and effectively of the uncrosslinked prepolymer and other residues, using a suitable solvent and can be blown dried with air. Disposable plastic molds inherently have variations in the dimensions, because, during injection-molding of plastic molds, fluctuations in the dimensions of molds can occur as a result of fluctuations in the production process (temperatures, pressures, material properties), and also because the resultant molds may undergo non-uniformly shrinking after the injection molding. These dimensional changes in the mold may lead to fluctuations in the parameters of contact lenses to be produced (peak refractive index, diameter, basic curve, central thickness etc.) and to a low fidelity in duplicating complex lens design. By using reusable molds which are produced in high precision, one can eliminate dimensional variations inherently presented in disposable molds and thereby variation in contact lenses produced therefrom. Lenses produced according to the Lightstream Technology™ can have high consistency and high fidelity to the original lens design.\nHowever, there are some practical limitations which hinder realization of all of the great potentials of such technology in the production of silicone hydrogel contact lenses. For example, when a silicone-containing prepolymer disclosed in commonly-owned U.S. Pat. Nos. 7,091,283, 7,268,189 and 7,238,750 is used to prepare a silicone hydrogel lens formulation, an organic solvent is generally required to solubilize the prepolymer. When such lens formulation is used to produce silicone hydrogel according to the Lightstream Technology™, the cured lens in the mold after UV crosslinking is still swollen in the organic solvent before the solvent exchange to water. Such lens may not be able to survive the mold opening and de-molding process, because the cured lens is in the swellon state in the organic solvent and has an inadequate stiffness and toughness (e.g., too low). As such, the production yield may be low and the production cost could be higher due to low production yield derived from the lens defects created during mold opening and de-molding process.\nAccordingly, there is still a need for a lens manufacturing process in which lens defects generated during mold opening and de-molding process can be minimized."}
{"text": "1. Field of the Invention\nThe present invention relates to a game processing server apparatus and a game processing server system.\n2. Description of the Related Art\nA social game is provided in a social networking service (SNS), and is an online game in which a player plays the game while having communications with other players. As the social game is provided in the SNS, which is proposed for providing communications between participants, the social game is configured to provide more communications to the players compared with a previously known online game although such an online game also provides communications to the players.\nAs one kind of such social games, a social Role Playing Game (RPG) is known. In the social RPG, a quest (mission) is provided in which a player character operated by a player (user) goes through a predetermined area in a map while playing the game.\nIn this quest, various events are set to occur when the player character goes through the predetermined area while consuming energy points or the like. For example, the player character can learn magic power, a skill or the like by encountering and defeating an enemy character or encountering a bonus item. Further, when the player clears a quest in an area where a boss character is set, the boss character appears so that the player character can fight against the boss character. When the player character defeats the boss character, the player character can get a weapon or a reward.\nThus, the player tries to clear quests and expects to obtain an ability, a weapon, a reward or the like through the predetermined events.\nIn a previously known RPG or in a social RPG, the event to occur in the same quest is set to be almost the same for all of the players based on a scenario of the game.\nEven in such a case, as long as the player enjoys the game by himself or herself, there may be no problem. Even when the same event occurs in the same quest for all of the players, the player recognizes the game as provided.\nHowever, for the social game, the players often have communications with other players, and sometimes, the players make a team to compete with other teams. Thus, the players often know what kind of event is to occur and what will happen as a result in a specific quest before the player actually plays the quest. In such a case, if the event occurs as the player expects, the player can easily predict the result of the game so that the player may have a tedious feeling thereby lowering the enjoyment of the game."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a charger, and more particularly to an AC and DC dual input charger.\n2. The Related Art\nA conventional charger is generally an AC power charger or a DC power charger. The AC power charger is directly connected with the AC power, and is mainly used inside a room. The DC power charger is connected with the DC power through a car charger, and is mainly used in a car and other vehicles. In order to make the charger used in both AC and DC work modes, an AC input plug and a DC input plug are provided to connect with the charger, and a manual control switch is mounted to the charger so as to switch the AC and DC input. However, when the AC and DC dual input charger described above is in use, it needs to control the switch manually to match the AC plug or DC plug so as to realize the function of charging. As a result, the usage of the AC and DC dual input charger becomes complex and apt to cause a wrong operation."}
{"text": "1. Field of the Invention\nThe present invention relates to a mounting structure with a heat sink for electronic component and securing members for the mounting structure and more particularly to the mounting structure with the heat sink for at least one electronic component suitably used for electronic components such as an LSI (Large Scale Integration) circuit whose temperature rises at a time of operations and the securing member suitably used for the mounting structure.\n2. Description of the Related Art\nConventionally, a heat sink of this kind is used for suppressing a temperature rise caused by power consumed when an LSI in a package is operating by being in thermal contact with the LSI surface-mounted or insertion-mounted on a printed circuit board. When the heat sink is secured to the printed circuit board, since the LSI package is sandwiched between the heat sink and the printed circuit board, it is necessary that some contrivance to accommodate variations in a height of the LSI package is provided.\nVariations in the height readily occur among the mounted LSI packages and reasons for the variations in the height of the LSI package are various. For example, in the case of a surface-mounting type LSI package such as a BGA (Ball Grid Array) or a like, variations in thickness of solder balls arranged in a grid-like form readily occurs among lots and, in the case of face-down bonding, variations occur easily, or also in an insertion-mounted type package, variations occurs readily in a degree to which a lead is inserted into a hole of a land portion.\nA related mounting structure of this type is disclosed in, for example, Patent Reference 1 (Japanese Patent Application Laid-open No. 2001-057405) in which a heat transmission rubber is disposed between a heat sink and an LSI package so that variations in a direction of height are accommodated by the heat transmission rubber.\nAlso, a related mounting structure of this kind is disclosed in, for example, Patent Reference 2 (Japanese Patent Application Laid-open No. Hei 09-139450) in which each compression coil spring is attached to each male screw serving as a securing member and is made to accommodate variations in height in a direction of an LSI package. However, thermal conductivity of the heat transmission rubber is low and highly-priced, which causes high costs of the heat sink securing structure, therefore, conventionally, the compression springs are generally used.\nFIG. 5 is an exploded perspective view showing a related mounting structure (structure for securing a heat sink) of an LSI package with a heat sink using compression coil springs and disassembled securing members to be used in the mounting structure. FIG. 6 is a side view showing the related mounting structure and the securing members for the mounting structure. FIG. 7 is the mounting structure of FIG. 5 taken along the line A-A.\nThe related mounting structure (structure for securing heat sink) of this type roughly includes, as shown in FIGS. 5, 6, and 7, a printed circuit board 1, an LSI package 2 being surface-mounted on the printed circuit board 1, a heat sink 3 having a fin structure disposed in a thermal contact state on an upper surface of the LSI package 2, a pair of lower and upper stiffeners (lower stiffener 4 and upper stiffener 5) to support and reinforce the printed circuit board 1, and a set of securing members including screw members 6, male screws 7, shafts 8, washers and a like, and compression coil springs 10 to accommodate variations in height in a direction of the LSI package 2 (refer to Patent Reference 3 [Japanese Patent Application Laid-open No. 2000-058703] for the stiffener structure). The pair of upper and lower stiffeners 4 and 5 is made of a metal plate and is secured with the screw members 6 in a state in which the printed circuit board 1 is sandwiched between the upper and lower stiffeners. The upper and lower stiffeners are secured to each other and, therefore, an occurrence of warpage caused by heat of the printed circuit board 1 can be prevented.\nThe above heat sink is so configured that many dissipating fins 3b to increase contact areas (dissipation area) with an outer atmosphere are disposed on an upper surface of an heat sink base 3a in parallel to one another and in an erected manner. At four corners of the heat sink base 3a are formed through-holes 3c to insert the male screws 7. On the upper surface stiffeners 6, as shown in FIG. 7, is formed an aperture portion to allow the LSI package to be inserted. In positions corresponding to the through-holes 3c at four corners of the heat sink base 3a are attached, in a securing manner, the shafts 8 (female screws with upper portion being opened) to screw the male screws 7 therein. The LSI package 2 is so configured as to be exposed from the aperture portion 11 of the upper stiffener 5 and a flat upper surface of the LSI package 2 is in thermal contact with a bottom of the heat sink base 3a. \nEach of the above male screws 7 is passed through each of washers 9, each of compression coil springs 10 and each of through holes 3c in this order and, in this state, an end of each of the male screws is screwed into each of the upper stiffener 5 and, as a result, each of the male screws 7 is screw-secured to each of the shafts. Thus, conventionally, the heat sink 3 is screw-secured to the upper stiffener 5 using the screw members and the printed circuit board 1, in a state in which the LSI package is exposed from the aperture portion 11 of the upper stiffener 5 is sandwiched between the upper stiffener 5 and the lower stiffener 4 and is secured and, therefore, in a state in which the upper surface of the LSI package 2 is in thermal contact with the bottom of the heat sink base 3a, the heat sink 3 is secured to the printed circuit board 1.\nAccording to the above configurations, even when variations in height occurs among the LSI packages 2 due to the easy occurrence of variations in the thickness of solder balls 12 (FIG. 7) arranged in a grid form and/or the easy occurrence of variations at a time of face-down bonding, each of the compression coil springs 10 sandwiched between a head portion of each of the male screws 7 and the upper surface of the heat sink base 3a, thereby variations in height can be accommodated (FIGS. 6 and 7).\nHowever, the above related mounting structure has a problem. That is, the height accommodating tool to secure the heat sink 3 to the printed circuit board 1 is made up of the male screws 7, the shafts (female screws) 5, the compression coil springs 10, and the washers 9 and, therefore, component counts are large, many attaching man-hours are required, thus causing complicated mounting processes. More specifically, the related method for the heat sink requires, as shown in FIG. 5, (1) a process of attaching and securing each of the shafts (female screws) 8 to the upper stiffener 5, (2) a process of securing the upper stiffeners 5 and the lower stiffener 4 by using each of the screw members 6 with the printed circuit board 1 being sandwiched between the upper and lower stiffeners, (3) a process of attaching each of the washers 9 and each of the compression coil springs 10 to each of the male screws 7, and (4) a process of securing the heat sink 3 to the printed circuit board 1 by using each of the male screws 7 with the heat sink 3 being in thermal contact with the LSI package 2, thus resulting in complicated mounting work. Additionally, the related mounting structure for the heat sink has another problem. That is, each of the male screws 7 and/or compression coil springs 10 spring out therefrom, a wind path F (FIG. 6) of a fan (not shown) is stopped up, thus resulting in lowering of the dissipation efficiency."}
{"text": "There are may characteristics of tissue products such as bath and facial tissue that must be considered in producing a final product having desirable attributes that make it suitable and preferred for the product's intended purpose. Improved softness of the product has long been one major objective, and this has been a particularly significant factor for the success of premium products. In general, the major components of softness include stiffness and bulk (density), with lower stiffness and higher bulk (lower density) generally improving perceived softness.\nWhile enhanced softness is a desire for all types of tissue products, it has been especially challenging to achieve softness improvements in uncreped throughdried sheets. Throughdrying provides a relatively noncompressive method of removing water from a web by passing hot air through the web until it is dry. More specifically, a wet-laid web is transferred from the forming fabric to a coarse, highly permeable throughdrying fabric and retained on the throughdrying fabric until dry. The resulting dried web is softer and bulkier than a conventionally-dried uncreped sheet because fewer bonds are formed and because the web is less compressed. Thus, there are benefits to eliminating the Yankee dryer and making an uncreped throughdried product. Uncreped throughdried sheets are typically quite harsh and rough to the touch, however, compared to their creped counterparts. This is partially due to the inherently high stiffness and strength of an uncreped sheet, but is also due in part to the coarseness of the throughdrying fabric onto which the wet web is conformed and dried.\nTherefore, what is lacking and needed in the art is a method for manufacturing tissue products having improved softness, and in particular uncreped throughdried tissue products having improved softness, as well as an apparatus that permits the manufacture of such tissue products."}
{"text": "The present invention relates to thermal barrier coatings, a method, and an apparatus for determination of past-service conditions of coatings and parts and remaining life thereof. In particular, the present invention relates to such a method and an apparatus by a non-destructive optical determination of a particular crystalline phase in a thermal barrier coating.\nThe constant demand for increased operating temperature in gas turbine engines has necessitated the development of ceramic coating materials that can insulate the turbine components such as turbine blades and vanes from the heat contained in the gas discharged from the combustion chamber for extending the life of such components. These ceramic coatings are known in the art as thermal barrier coatings.\nA thermal barrier coating typically comprises at least a layer of a refractory or thermally insulating material such as yttria-stabilized zirconia (or xe2x80x9cYSZxe2x80x9d) which is zirconia stabilized with, for example, about 6-8 percent by weight of yttria. The refractory material would generally be selected to have a low thermal conductivity such as about 1-3 W/(m)(K), thereby reducing heat transfer to and the temperature experienced by the turbine engine component. The coating may be applied by one of known deposition techniques such as the thermal or plasma spray process or the physical vapor deposition process. A typical thermal barrier coating is a multilayer system comprising three layers. A first so-called bondcoat is applied to the surface of the superalloy of the turbine component. This bondcoat typically comprises a MCrAlY alloy wherein M is nickel, or cobalt, or PtNiAl alloys. The purpose of the bondcoat is to provide a layer which adheres well to the underlying alloy, which provides protection against oxidation of the alloy, and which provides a good base for further coatings. A second intermediate layer or interlayer is applied on the bondcoat. A suitable material for this interlayer is Al2O3. This material can be formed by oxidizing the surface of the bondcoat to form an oxide layer. The interlayer provides improved adhesion for the final thermal insulating YSZ coating and is not included for a thermal barrier property.\nDespite great care taken during manufacture to ensure good adhesion of the thermal barrier coating to the underlying material of the turbine component, thermal cycling during use of such a component eventually leads to spalling of the coating. In addition, erosion of the thermal barrier coating is inevitable over an extended period of use. Such a spalling or erosion would eventually expose the underlying alloy to extreme temperatures that would lead to failure of the component. Therefore, thermal barrier coatings need be inspected frequently for any sign of deterioration. Such an inspection often requires taking the engine component out of service and is time-consuming. A common inspection technique is the visual inspection of the presence or absence of coating. While that method determines when a spall has occurred, it is unable to determine either the degree of deterioration in an intact coating. A method for determining the past-service conditions and remaining life of thermal barrier coatings would be welcome in the art.\nSimilarly, it is desirable to monitor the condition of the turbine components themselves. In the prior art, it is usual for a destructive evaluation to be performed at each inspection interval for critical components in the hot gas path. In that case, one part is destroyed to produce sections for metallographical examination. The condition of the coatings and base materials are determined from metallographical inspection, and a decision to repair or replace the remaining parts is made from that information.\nBetter knowledge of the past-service conditions experienced by the turbine components would allow the determination of the remaining life of a part without destructive evaluation. Currently there are few in-situ measurements of hot gas path parts temperatures available. Some physical changes in the phases and structures of the materials of thermal barrier coatings and components occur with exposure to high temperatures. Inspection for changes in phase content is one way to determine past-service conditions.\nHowever, traditional methods of inspection, such as X-ray diffraction and neutron diffraction, require destructive testing and specialized equipment. They are not conducive to being deployed at the site of a gas turbine. In addition, such destructive testing methods necessarily extrapolate the result obtained for one part to the condition of other similarly used parts and, thus, may not provide a true and accurate condition of those parts.\nEuropean patent application EP 0863396 A2 discloses a non-destructive measurement method for residual stress proximate an interlayer in a multilayer thermal barrier coating system. This method focuses on detecting compressive stresses that accumulate at the boundary between the interlayer and the outermost thermal barrier coating by detecting the shift in frequency of light emitted by fluorescing chromium ions in the alumina interlayer. However, significant stresses at that boundary may not appear until after the outermost barrier layer has seriously deteriorated. Furthermore, the stresses at the boundary are not useful indicators of the past-service conditions of the component itself. Therefore, such a method is not very useful in timely forewarning a need for repairing or replacing the engine component.\nTherefore, there is a continued need to provide a simple non-destructive method for determining the past-service condition of a thermal barrier coating of a component used at high temperature in a turbine engine. It is also very desirable to provide a method by which the remaining useful life of the underlying component may be determined or estimated. Furthermore, it is also very desirable to provide such a method so that maintenance of turbine engine components may be performed only on an as-needed basis rather than on a fixed schedule.\nThe present invention provides a method for determining at least one of past-service conditions and remaining useful life of at least one of a component of a combustion engine and a thermal barrier coating thereof, which component is used in the hot-gas path of the combustion engine. The method of the present invention comprises (1) providing a combustion-engine component comprising a thermal barrier coating that comprises at least one photoluminescent (xe2x80x9cPLxe2x80x9d) material that can be excited by radiation at a first wavelength range and emits radiation at a second wavelength range different from the first wavelength range in response to the exciting radiation; the radiation emitted at the second wavelength range having a characteristic property that correlates with an amount of a crystalline phase in the thermal barrier coating, which amount increases as the combustion-engine component is exposed to elevated temperatures; (2) directing radiation having the first wavelength range at the thermal barrier coating of the combustion-engine component; (3) measuring the characteristic property of radiation having the second wavelength range; (4) determining the amount of the crystalline phase present in the thermal barrier coating from the characteristic property of radiation having the second wavelength range; and (5) inferring at least one of past-service conditions and remaining useful life of the thermal barrier coating from the amount of the crystalline phase.\nAccording to one aspect of the present invention, the thermal barrier coating comprises yttria-stabilized zirconia.\nAccording to another aspect of the present invention the crystalline phase is the monoclinic phase of zirconia.\nThe present invention also provides an apparatus for determining at least one of past-service conditions and remaining useful life of at least one of a component of a combustion engine and a thermal barrier coating thereof, which component is used in the hot-gas path of the combustion engine. The apparatus comprises (1) a source of radiation having a first wavelength range directed at the thermal barrier coating that comprises at least one PL material capable of emitting radiation having a second wavelength range in response to an excitation by the radiation having the first wavelength range; (2) a radiation detector being capable of detecting the radiation having the second wavelength range and being disposed to receive and measure a characteristic property thereof; and (3) means for relating the characteristic property of radiation having said second wavelength range to one of an amount of a crystalline phase, past-service conditions, and remaining useful life of the combustion-engine component.\nOther features and advantages of the present invention will be apparent from a perusal of the following detailed description of the invention and the accompanying drawings in which the same numerals refer to like elements."}
{"text": "Every clothes dryer requires an exhaust channel to function efficiently. Normally, the exhaust channel has a hose attached to a dryer connector at one end and wall or floor vent at the other. In order to simplify the attaching of the exhaust channel, a flexible hose is known and used. However, for safety reasons, in some cases, use of the flexible hose is restricted. With such a restriction, installation of a clothes dryer becomes more difficult.\nA clothes dryer is difficult to move. Yet, movement of the clothes dryer is mandatory for connection and cleaning of the exhaust channel. If a method can be found, which reduces the movement of the clothes dryer while installing the exhaust channel, great advantages can be obtained.\nIf the flexible hose is not used, a metal pipe is the desired replacement for the exhaust channel. Since the installation of the dryer is often in a very confined space, it becomes difficult to maneuver and attach the metal pipe to the dryer at one end and the exhaust vent at the other end, so a minimal length of flexible hose is used to connect an exhaust channel to a dryer exhaust outlet. While the flexible hose stands a better chance of collecting lint than the metal pipe, both types do collect lint and thereby present a fire hazard. As the length of the exhaust hose or pipe increases, the danger of lint collection increases, as does the fire danger. In order to keep the exhaust channel free of lint during operation of the dryer, the exhaust pipe must be detached from the clothes dryer occasionally to give access for inspection and removal of any accumulated lint. This is very desirable to minimize the danger of such fire. Thus, the prior art suffers from numerous drawbacks for the connection and disconnection of the exhaust channel. If the connection and disconnection can be simplified, while improving fire safety, great advantages can be obtained."}
{"text": "Since the late 1980's hazardous agents, such as cytotoxic agents have been useful in managing and treating a number of diseases such as rheumatoid arthritis (and other autoimmune diseases), juvenile rheumatoid arthritis, psoriatic arthritis, systemic lupus erythematosus, steroid resistant polymyositis or dermatomyositis, Wegener's granulomatosis, polyarteritis nodosa, and some forms of vasculitis. Hazardous agents tend to exhibit side effects, however, that are harmful or toxic to the subject. Many of these side effects occur when hazardous agents are administered orally, but the oral form is generally the preferred method of delivery of these agents due to its ease of use.\nIn addition to increased toxicity, variable and reduced bioavailability has been observed for some hazardous agents, such as methotrexate, that are orally administered. These limitations are particularly demonstrated when the oral dosing is escalated beyond 15 mg per week. It has been suggested that with parenteral administration, such as by injection, more predictable, reproducible and complete bioavailability along with better therapeutic results could be achieved, particularly at higher dosages.\nOnly about 7% of the prescriptions for methotrexate written by rheumatologists are for an injectable formulation. Reasons for prescribing methotrexate injections are usually to improve bioavailability or to alleviate side effects. Physicians have expressed interest in increasing the number of prescriptions for cytotoxic agent injections, and particularly injections for home use and administration by a patient. This is generally not considered feasible because it is not possible to ensure that patients can reliably and repeatably draw an accurate dose from vials and correctly administer the product by subcutaneous (SC) injection, especially with agents used to treat patients suffering from certain debilitating diseases. Additionally, the toxicity of hazardous agents increases the risk that non-users of the injections will come into contact with the cytotoxic agents in a home setting. Insufficient data exists on the effect of low dose, chronic exposure to hazardous agents that are, or may be, candidates for home use or self-injection. In the absence of such information, practice guidelines direct one to assume a high degree risk for injectable hazardous agents such as methotrexate, with the recommendation of formal directives and risk assessments, including formal training and mitigation strategies, to minimize risk (see Oliver, S., and Livermore, P., Administering subcutaneous methotrexate for inflammatory arthritis: RCN guidance for nurses, 2004; Royal College of Nursing, Wyeth, Publication Code 002 269). Specific directives include: preparation of syringes in dedicated pharmacies with aseptic preparation areas; administration performed in specific locations and only by adequately trained personnel; spillage kits located proximal to use areas; accounting for all who may be at risk in the event of an accident; and audits to assess compliance and execution of risk mitigation strategies. Because of the need for such directives, and thus the large number of precautions that must be learned and followed in order to safely inject a hazardous agent, it is presently thought that it is not practical for hazardous agents, and particularly methotrexate, to be self-injected by a patient outside of a clinical setting or without the assistance of a health care provider."}
{"text": "1. Field of the Invention\nThis invention is directed to a unique method and device for delivering controlled heat to perform ablation to treat benign prosthetic hypertrophy or hyperplasia (BPH). The method and the apparatus deliver this controlled heat into tissue penetrated by devices such as those disclosed in the copending above-referenced applications.\n2. Discussion of Background\nTreatment of cellular tissues usually requires direct contact of target tissue with a medical instrument, usually by surgical procedures exposing both the target and intervening tissue to substantial trauma. Often, precise placement of a treatment probe is difficult because of the location of a target tissue in the body or the proximity of the target tissue to easily damaged, critical body organs, nerves, or other components.\nBenign prostatic hypertrophy or hyperplasia (BPH), for example, is one of the most common medical problems experienced by men over 50 years old. Urinary tract obstruction due to prostatic hyperplasia has been recognized since the earliest days of medicine. Hyperplastic enlargement of the prostate gland often leads to compression of the urethra, resulting in obstruction of the urinary tract and the subsequent development of symptoms including frequent urination, decrease in urinary flow, nocturia, pain, discomfort, and dribbling. The association of BPH with aging has been shown to exceed 50% in men over 50 years of age and increases in incidence to over 75% in men over 80 years of age. Symptoms of urinary obstruction occur most frequently between the ages of 65 and 70 when approximately 65% of men in this age group have prostatic enlargement.\nCurrently there is no proven effective nonsurgical method of treatment of BPH. In addition, the surgical procedures available are not totally satisfactory. Currently patients suffering from the obstructive symptoms of this disease are provided with few options: continue to cope with the symptoms (i.e., conservative management), submit to drug therapy at early stages, or submit to surgical intervention. More than 30,000 patients per year undergo surgery for removal of prostatic tissue in the United States. These represent less than five percent of men exhibiting clinical significant symptoms.\nThose suffering from BPH are often elderly men, many with additional health problems which increase the risk of surgical procedures. Surgical procedures for the removal of prostatic tissue are associated with a number of hazards including anesthesia associated morbidity, hemorrhage, coagulopathies, pulmonary emboli and electrolyte imbalances. These procedures performed currently can also lead to cardiac complications, bladder perforation, incontinence, infection, urethral or bladder neck stricture, retention of prostatic chips, retrograde ejaculation, and infertility. Due to the extensive invasive nature of the current treatment options for obstructive uropathy, the majority of patients delay definitive treatment of their condition. This circumstance can lead to serious damage to structures secondary to the obstructive lesion in the prostate (bladder hypertrophy, hydronephrosis, dilation of the kidney pelves, etc.) which is not without significant consequences. In addition, a significant number of patients with symptoms sufficiently severe to warrant surgical intervention are poor operative risks and are poor candidates for prostatectomy. In addition, younger men suffering from BPH who do not desire to risk complications such as infertility are often forced to avoid surgical intervention. Thus the need, importance and value of improved surgical and non-surgical methods for treating BPH is unquestionable.\nHigh-frequency currents are used in electrocautery procedures for cutting human tissue especially when a bloodless incision is desired or when the operating site is not accessible with a normal scalpel but presents an access for a thin instrument through natural body openings such as the esophagus, intestines or urethra. Examples include the removal of prostatic adenomas, bladder tumors or intestinal polyps. In such cases, the high-frequency current is fed by a surgical probe into the tissue to be cut. The resulting dissipated heat causes boiling and vaporization of the cell fluid at this point, whereupon the cell walls rupture and the tissue is separated.\nDestruction of cellular tissues in situ has been used in the treatment of many diseases and medical conditions alone or as an adjunct to surgical removal procedures. It is often less traumatic than surgical procedures and may be the only alternative where other procedures are unsafe. Ablative treatment devices have the advantage of using a destructive energy which is rapidly dissipated and reduced to a non-destructive level by conduction and convection forces of circulating fluids and other natural body processes.\nMicrowave, radiofrequency, acoustical (ultrasound) and high energy (laser) devices, and tissue destructive substances have been used to destroy malignant, benign and other types of cells and tissues from a wide variety of anatomic sites and organs. Tissues treated include isolated carcinoma masses and, more specifically, organs such as the prostate, glandular and stromal nodules characteristic of benign prostate hyperplasia. These devices typically include a catheter or cannula which is used to carry a radiofrequency electrode or microwave antenna through a duct to the zone of treatment and apply energy diffusely through the duct wall into the surrounding tissue in all directions. Severe trauma is often sustained by the duct wall during this cellular destruction process, and some devices combine cooling systems with microwave antennas to reduce trauma to the ductal wall. For treating the prostate with these devices, for example, heat energy is delivered through the walls of the urethra into the surrounding prostate cells in an effort to kill the tissue constricting the urethra. Light energy, typically from a laser, is delivered to prostate tissue target sites by \"burning through\" the wall of the urethra. Healthy cells of the duct wall and healthy tissue between the nodules and duct wall are also indiscriminately destroyed in the process and can cause unnecessary loss of some prostate function. Furthermore, the added cooling function of some microwave devices complicates the apparatus and requires that the device be sufficiently large to accommodate this cooling system.\nApplication of liquids to specific tissues for medical purposes is limited by the ability to obtain delivery without traumatizing intervening tissue and to effect a delivery limited to the specific target tissue. Localized chemotherapy, drug infusions, collagen injections, or injections of agents which are then activated by light, heat or chemicals would be greatly facilitated by a device which could conveniently and precisely place a fluid supply catheter opening at the specific target tissue."}
{"text": "The present invention relates to a buffer storage system in a data processing apparatus and, more particularly, it relates to a buffer storage system including therein both a buffer storage for accessing operands and a buffer storage for fetching instructions, separately.\nRecent increases, in the capability of a data processing apparatus, have led to attempted use of a variety of data-processing methods for coping with this increased capability. In one of these methods, namely the pipeline method, each sequence for executing an instruction is divided into a plurality of phases and each of the phases is executed at each station. This enables operations at respective stations to be achieved independently, and a plurality of instructions to be executed simultaneously.\nIn the prior art buffer storage system, both the pipeline set up in the operand access buffer storage, and the pipeline set up in the instruction fetch buffer storage, must be synchronized with each other. To strictly maintain such synchronization, often several operand access cycles must be left idle until several instruction fetch cycles with high priority are completed. This will be clarified in detail hereinafter. However, this leads to the problem in that the effective operating speed of the data processing apparatus cannot be made sufficiently high, due to the insertion of the above-mentioned idle operand access cycles, and therefore, the capability of the data processing apparatus cannot be increased to any large degree, even though independent buffer storages for accessing operands and fetching instructions are introduced."}
{"text": "1. Field of the Invention\nThe invention relates to a method of manufacturing a semiconductor device, and particularly to a method of manufacturing a non-volatile memory.\n2. Description of Related Art\nWhen semiconductor technology enters deep sub-micron manufacturing process, sizes of devices are gradually decreased, which means decreased memory cell size with respect to memory device. On the other hand, as data which information electronic products (such as computer, mobile phone, digital camera and personal digital assistant (PDA)) have to handle and store is increasing, the memory capacity required by these information electronic products becomes larger and larger. In the case of decreased device size and increased memory capacity demand, a common goal in the field is how to manufacture memory devices having decreased size and high integration while maintaining good qualities.\nA non-volatile memory is capable of safeguarding stored data even after the power supplied to the non-volatile memory is cut off, and therefore the non-volatile memory has been extensively applied to personal computers and electronic equipments.\nA typical non-volatile memory cell has a memory gate and a control gate made by doped polysilicon. A dielectric layer is disposed respectively between the memory gate and a substrate, and between the control gate and the substrate.\nHowever, the manufacture of the above-mentioned non-volatile memory cell requires forming a plurality of polysilicon layers and a plurality of dielectric layers. During the manufacturing process, several photomasking steps are carried out, which not only lengthens the manufacturing process but also incurs more manufacturing cost.\nA conventional NOR type non-volatile memory cell formed by two transistors connected in series includes a select transistor and a memory gate transistor. For this type of memory cell, there is no need to form a plurality of polysilicon layers. Hence the manufacturing process of such non-volatile memory cell can be integrated with the manufacturing process of a complementary metal oxide semiconductor transistor.\nGenerally, the non-volatile memory consists of a plurality of memory cells located in a memory cell region and a plurality of logic devices (such as input/output transistor, core transistor, etc.) located in a periphery circuit region. The select transistor and the input/output transistor in the periphery circuit region are manufactured in the same process. A gate dielectric layer of the input/output transistor is usually thicker for withstanding a higher operating voltage. However, in a situation where the sizes of devices are decreased due to increased integration of an integrated circuit, the size of the memory cell is decreased as well. If a thickness of a gate dielectric layer of the select transistor is equal to the thickness of the gate dielectric layer of the input/output transistor in the periphery circuit region, during operation of the memory, a larger voltage needs to be applied to a gate of the select transistor, with the result that a driving capability of a non-volatile memory device is reduced. Therefore, it will be an important issue that how to enable the non-volatile memory device to have a better driving capability."}
{"text": "The preparation of diazaspirodecan-2-ones, in particular, 8-[{1-(3,5-Bis-(trifluoromethyl)phenyl)-ethoxy}-methyl]-8-phenyl-1,7-diaza-spiro[4,5]decan-2-ones, for example, (5S,8S)-8-[{(1R)-1-(3,5-Bis-(trifluoromethyl)phenyl)-ethoxy}-methyl]-8-phenyl-1,7-diazaspiro[4,5]decan-2-one (the compound of Formula I) is disclosed in U.S. Pat. No. 7,049,320, issued May 23, 2006 (the '320 patent) which is incorporated herein by reference in its entirety.\n\nThe novel compounds disclosed in the '320 patent are classified as Tachykinin compounds, and are antagonists of neuropeptide neurokinin-1 receptors (referred to herein for convenience as “NK-1 receptor antagonists”).\nThe compounds described in the '320 patent are classified as tachykinin compounds, and are antagonists of neuropeptide neurokinin-1 receptors (herein, “NK-1” receptor antagonists). Other NK-1 receptor antagonists and their synthesis have been described, for example, those described in Wu et al, Tetrahedron 56, 3043-3051 (2000); Rombouts et al, Tetrahedron Letters 42, 7397-7399 (2001); Rogiers et al, Tetrahedron 57, 8971-8981 (2001) and in each of the following publications: published international application no. WO05/100358; U.S. Pat. No. 5,760,018 (1998); U.S. Pat. No. 5,620,989 (1997), and international publication nos. WO 95/19344 (1995), WO 94/13639 (1994), and WO 94/10165 (1994), each of which are incorporated herein in their entirety by reference.\n“NK-1” receptor antagonists have been shown to be useful therapeutic agents, for example, in the treatment of pain, inflammation, migraine, emesis (vomiting), and nociception. The novel NK-1 compounds disclosed in the above-mentioned '320 patent include the compound of Formula I, which is useful in the treatment of nausea and emesis associated with chemotherapy treatments (Chemotherapy-induced nausea and emesis, CINE). Emesis and nausea have been a problem in the provision of chemotherapy. Chemotherapeutic agents, for example, cisplatin carboplatin and temozolomide have been associated with both acute and delayed onset nausea and vomiting. It is known to administer chemotherapeutic agents with an anti-emetic, for example, as described in U.S. Pat. No. 5,939,098, which describes coadministration of temozolomide and with ondansetron, however such therapy is not effective in preventing delayed onset nausea and vomiting.\nAs reported in the '320 patent, the compound of Formula I was characterized by TLC and by GC/MS techniques. The procedures described in the '320 patent yielded the compound of Formula I in the form of an amorphous white foam. Repeated attempts to crystallize the free base have not provided a crystalline material.\nIn general, compounds which have been identified as having therapeutic activity must be provided in a highly pure form for pharmaceutical use. Moreover, it is desirable to provide compounds intended for pharmaceutical use in a form such that it is handled easily for incorporation into a medicament, and when incorporated into a medicament the compound possesses a sufficiently robust character that it is resistant to chemical degradation, and thereby imparts a long shelf life to the medicament."}
{"text": "Lithium ion secondary batteries have been used in recent years as a power supply with ability of charge and discharge of electricity which is incorporated in cellular phones or portable electric devices. Moreover, for example, batteries using a solid-like unfluidized electrolyte as the electrolyte without danger of liquid leakage are also known. There are various outside forms of such a lithium ion secondary battery, and the batteries, generally used for notebook type or pocketbook type portable electric devices and cellular phones, are flat type in many cases. In such a flat type lithium ion secondary battery, in order to produce electric current continuity between a main body of a secondary battery cell and the exterior, board-like electrodes are respectively located on the positive pole/terminal and the negative pole/terminal.\nThe pole/terminal of positive electrode is generally formed with thin aluminum or an aluminum base alloy having a thickness of about 0.07-0.1 mm by the press cut or the like, in order to satisfy the structural limitations in a part connected to the main body of the secondary battery cell, restrain and reduce the thickness of the whole secondary battery, and ensure high conductivity, and a metal plate (the so-called tab) further connected with the exterior terminal lead is welded to the tip portion in many cases. As the metal plate, materials with low electric resistance, excellent mechanical strength, and excellent weather resistance such as a nickel base alloy are used suitably in order to electrically and mechanically (based on strength of material) ensure the connection with the exterior. Moreover, the plate thickness is set, for example, to 0.1 mm or more in many cases.\nThe metal plate and the electrode are generally welded by the electric resistance welding or the ultrasonic welding, and the sure welding is strongly needed so that the metal plate may not be separated nor omitted from the electrode, while using the secondary battery or the like.\nHowever, there is a problem that poor welding occurs in the case of the electric resistance welding between the above electrode and the metal plate.\nThe first reason for the difficulty of the electric resistance welding is the difference between aluminum and nickel melting points. That is, the pole/terminal of positive electrode is formed with aluminum or the aluminum base alloy, so the melting point is 660-700° C. On the other hand, the metal plate is formed with the material with both high mechanical strength and a comparatively high melting point, so the melting point is, for example, 1400-1455° C. in the case of the nickel base alloy. Therefore, the difference between the electrode and the metal plate melting points reaches about 800° C. Moreover, the boiling point of aluminum is 2486° C., the boiling point of nickel is 2731° C., and the boiling points also differ greatly.\nThe second reason for the difficulty of the electric resistance welding is the existence of an oxide film (aluminum oxide) formed on the surface of the aluminum plate. The melting point of aluminum oxide is as high as 2050° C., and in the case of the electric resistance welding between the aluminum plate and the nickel plate, it is necessary to dissolve the thin oxide film under the temperature of the weld part of about 2050° C. or higher. Here, the aluminum oxide film on the aluminum surface is generally called alumina, and the chemical formula thereof is Al2O3.\nThus, if a welding condition is set up so that the welding temperature may reach the melting point or more of the metal plate made of the nickel base alloy or the like, in the case of welding a pile of one electrode and one metal plate by the conventional general electric resistance welding process, aluminum or the aluminum base alloy of the electrode dissolves completely over a large area to the extent that the plate thickness is penetrated exceeding the size of a normal/regular nugget. Accordingly, the dissolved metal is spilt out, aluminum evaporates and scatters around violently after the aluminum plate of the weld part reaches the boiling point, or the like phenomenon occurs, which causes generation of a hole in the part and poor welding without the sure welding.\nMoreover, when the welding temperature is adjusted under the melting point of the metal plate, such as the nickel base alloy, in order to avoid such a complete broad dissolution of the electrode, the poor welding such as the hole generation in the electrode or the spill of the dissolved metal does not occur. However, since the nickel base alloy on the metal plate surface does not fuse, the normal nugget of both the nickel base alloy and the aluminum base alloy is not formed, only the trace of the aluminum base alloy melting is left on the cross section surface after the welding, and the poor welding of the unsure welding occurs. Consequently, for example, the peeling test reveals that the electrode and the metal plate separate easily under very lower power than the specific tensile proof stress.\nThen, it is desired that various conditions of the electric resistance welding are set in order to generate a distribution of the welding temperature in which the normal nugget can be formed from the junction surface between the metal plate and the electrode to the peripheral part. However, such a condition setup unavoidably becomes very delicate, since the difference between the metal of the metal plate such as the nickel base alloy and metal of the electrode such as the aluminum base alloy melting points, is too large. Moreover, the tip of the electrode pole is degraded and deformed as the welding is continued, so it is very difficult to continuously maintain the setup of the above preferable delicate welding conditions in the actual mass production process, and therefore it is difficult to reduce the occurrence rate of the poor welding.\nIn addition, if the positive electrode terminal strip is composed of a nickel plate like the negative electrode terminal strip, the electric resistance welding can be carried out. However, if the nickel plate composes the positive electrode terminal strip, the nickel plate may react electrically in the inside of the battery and dissolve, and thus it is not preferable.\nAs described above, it is technically difficult to perform the electric resistance welding of the thin aluminum plate and the nickel plate, and particularly, the welding of the positive electrode terminal strip and a wiring board is conventionally performed by the ultrasonic welding.\nHowever, there have been the following problems in the ultrasonic welding. First, the setting ranges of oscillating strength and amplitude time of the ultrasonic welding are small, and it is difficult to maintain the optimum conditions of the welding. Secondly, it is difficult to stabilize the welding strength, and the poor welding may occur at a certain rate in the mass production process for manufacturing the products in large quantities. Because, in the ultrasonic welding, only a very thin alloy layer is produced at the interface between the electrode and the metal plate, both are in the weakly connecting state only on the surface in many cases due to increasing the roughness of both the surfaces, and it is difficult to accomplish the demanded electrically and mechanically sure welding state for the welded surface. Moreover, the natural oxide film generated on the surface of the metal plate made of aluminum or the aluminum base alloy presents obstacles, and the weak welding action by the ultrasonic welding tends to be still weaker. And if the electrode and the metal plate are welded weakly as described above, the electric resistance becomes high in the part thereof, with such unfortunate consequences as the voltage which can be transferred from the secondary battery cell outside through the electrode and the metal plate decreases or the electric current can be limited.\nThirdly, ultrasonic welding devices are more expensive than resistance welding devices, and the equipment expenses for mass production becomes high. Fourthly, the ultrasonic welding devices are larger than the resistance welding devices, and a larger floor space is required. Then, developments in the technique for welding the thin aluminum plate and the nickel plate by the electric resistance welding have been demanded.\nIt is necessary, for example, to weld for long time using the supersonic wave with still stronger energy density or the like in order to produce the still stronger welding state by such ultrasonic welding. However, if the welding is performed with such strong supersonic wave for long time, the electrode, made of the aluminum base alloy with the thin thickness and the lower melting point, dissolves in the range exceeding the size required for welding like the case of the above electric resistance welding, and the advantages of the ultrasonic welding itself cannot be employed efficiently. For example, if too strong ultrasonic welding is performed, a part of the electrode made of the aluminum base alloy is cut.\nMoreover, although the utilization of the soldering method is also considered as a method of electrically and mechanically fixing the metal plate to the electrode other than the welding, the natural oxide film (alumina) as described above is generated on the surface of the electrode made of the aluminum base alloy, and therefore presents obstacles to the wettability and attachment of the solder, and the soldering becomes difficult. It is thought that the treatment of pre-applying powerful flux or the like to the surface of the electrode for processing the natural oxide film of the surface before the soldering is effective to overcome the above obstacles. However, the inevitable result is that the component of such a strong flux remains on the electrode or the metal plate after soldering, so there are unfortunate consequences that the remaining component may remarkably degrade the durability of the connection part between the electrode and the metal plate. For example, the remaining components may gradually corrode the electrode during the period of using the secondary battery, which would eventually be damaged, dropped out, or the like. Moreover, soldering the positive electrode terminal strip and the wiring board is not preferable, since the thin aluminum plate composing the positive electrode terminal strip is brought to high temperature, and the inside temperature of the battery also becomes high, resulting in the battery degradation.\nThe present invention has been achieved in view of the above problems. It is an object of the invention to provide a method of manufacturing a battery comprising the step of welding an electrode made of aluminum, an aluminum base alloy, or the like, and a metal plate made of a nickel base alloy or the like, in a polymer lithium ion secondary battery or the like, by means of an electric resistance welding process comprising an electric resistance welding step of securely welding the electrode and the metal plate, with eliminating the problems of poor welding such as hole generation or scattering to the surroundings due to the spilled metal in a welding part, and a battery with high reliability and durability in which a metal plate and an electrode are securely welded by such method.\nIt is another object of the invention to provide a method of manufacturing a weldment in which two or more objects being welded which made of a different material are easily welded by the electric resistance welding with high reliability, and a pedestal used for the method."}
{"text": "Field of the Invention and Related Art Statement\nThe present invention relates to an automatic analyzer comprising a reaction line along which a number of reaction vessels each containing respective test liquids are transported and photometering means for effecting the photometry for the test liquids in the reaction vessels by transmitting light beams through the reaction vessels.\nAn automatic analyzer of single-line and multi-item type has been proposed in which a plurality of test items, i.e. a plurality of substances in samples are analyzed by means of a single reaction line. In such an automatic analyzer, test liquids contained in reaction vessels have to be photometered with the aid of light beams having different wavelengths corresponding to respective test items.\nFIG. 1 is a schematic view showing a known photometering apparatus disclosed in Japanese Patent Publication No. 65-21,303. The photometering apparatus comprises a light source 1 emitting white light, i.e. light including whole wavelength components, a condenser lens 2 for collecting the white light and making it incident upon reaction vessels 4 transported along a reaction line 3. Light transmitted through a test liquid contained in a reaction vessel 4 is made incident upon a spectroscope 6 via a slit 5 and is divided into a plurality of light beams having different wavelengths. These light beams are then made incident upon a plurality of light detectors 8-1, 8-2, . . . 8-n, respectively by means of a slit 7 having a plurality of holes 7a. One or more output signals supplied from the light detectors 8-1, 8-2, . . . 8-n are selected in accordance with test items to be measured for respective samples.\nFIG. 2 is schematic view illustrating another known photometering apparatus described in Japanese Patent Laid-open Publication, Kokai Sho 60-117,118. In this known apparatus, white light emitted from a light source 11 is evenly made incident upon incident ends of a plurality of optical fibers 13 via a condenser lens 12. Exit ends of optical fibers 13 are secured at positions which are predetermined in accordance with wavelengths of measuring light beams emanating from a diffraction grating 14. In opposition to the exit ends of optical fibers 13 there is arranged a rotary disc 19 which is rotated by a motor 17. As shown in FIG. 3, the rotary disc 19 has formed therein a plurality of sector slits 18 corresponding to the positions at which the exit ends of optical fibers 13 are arranged. The motor 17 is driven such that any one of slits 18 in the disc 19 can be positioned opposite an exit end of an optical fiber in accordance with a test item destined for a test liquid contained in a reaction vessel 16 which is just indexed at a measuring position defined by slit 20. Therefore, a light beam having a desired wavelength is made incident upon the reaction vessel 16 from the diffraction grating 14, and the light transmitted through the reaction vessel is received by a light detector 21. The reaction vessels 16 are transported along a reaction line 15 through the measuring position.\nIn the known photometering apparatus illustrated in FIG. 1, since the white light has a large amount of energy and is made incident upon the test liquid, some substances in test liquids might be decomposed or altered, so that in practice it is difficult to carry out the measurement precisely.\nThis problem could be removed by the known apparatus illustrated in FIG. 2, because only a slight flux having a desired wavelength corresponding to a test item is made incident upon a test liquid. However, in this known apparatus, in order to select the wavelength it is necessary to rotate the slit disc 19, and therefore a long time period may be necessary for measuring test items, thereby reducing the processing ability of such systems. Moreover, the white light emitted from the light source 11 is divided into a plurality of light beams with the aid of the optical fibers and thus weaking the intensity of respective light beams. And as a result, noise can affect the weakened photometered signals thus decreasing the accuracy of the photometry.\nU.S. Pat. No. 4,528,159, issued on July 9, 1985, discloses another known automatic analyzer comprising a light source emitting white light, first and second filter wheels arranged rotatably, and first and second light guides arranged between a cuvette and the filter wheels, respectively. Light emitted from the light source is evenly made incident upon filter elements of the first filter wheel, and a light flux emanating from a filter element is made incident upon a cuvette containing a test liquid via the first light guide. A light flux transmitted through the cuvette is made incident upon a light detector by means of the second light guide and a filter element of the second filter wheel. This known photometering apparatus is principally the same as the known apparatus shown in FIG. 2, and thus it requires a rather long time period to rotate the first and second filter wheels so as to index desired filter elements thereof into the measuring optical path in accordance with a test item to be measured."}
{"text": "Commonly assigned U.S. Pat. No. 5,448,430 is incorporated for its showing of a track following servo system for following servo track edges of dissimilar servo signals, commonly assigned U.S. Pat. No. 5,844,814 is incorporated for its showing of an independent position sensor in a head positioning system, and commonly assigned U.S. Patent Application (Ser. No. 09/413,327) is incorporated for its showing of a servo position detector and a method for detecting and following an index servo position displaced with respect to an edge of a servo track.\nThis invention relates to recording system track following servos, and, more particularly, to the initialization and calibration of indexed servo positions displaced from servo track edges.\nAdvancements in technology in the data storage industry often include increases in the data storage capacity of given data storage media. One means of increasing the data storage capacity of data storage media, such as magnetic tape cartridges or magnetic tape cassettes, is to increase the track density of the data storage media.\nIn a typical magnetic tape, data is recorded in a plurality of parallel, longitudinal data tracks. A tape head may have a plurality of data heads which have fewer numbers of read/write elements than tracks. The data tracks are divided into groups, typically interleaved, and the tape head is indexed laterally with respect to the tracks to access each group of data tracks. In order to properly register the read/write elements with the data tracks, prerecorded servo tracks are provided which are parallel to the data tracks. A servo head located at an indexed position with respect to the read/write elements reads the servo tracks. The servo tracks provide lateral positioning information which, when read by the servo head, can be detected by a servo detector to indicate whether the servo read head is correctly positioned with respect to the servo tracks. Thus, the servo head can be moved laterally to a desired position with respect to the servo tracks so as to properly register the read/write elements with respect to a desired group of data tracks. Then, the servo head can follow the servo tracks as the media and the head are moved longitudinally with respect to each other, so that the read/write elements maintain registration with the data tracks. Typically, the servo head follows servo tracks at an edge, an edge comprising an interface between two dissimilar recorded servo signals.\nThe incorporated (Ser. No. 09/413,327) application utilizes existing servo tracks having edges, but increases the data track density by employing index servo positions displaced from the edges. Effectively, the index servo position is the lateral position on the tape at which the center of the servo read head is located, and this position is laterally displaced from an edge. The edges are easily tracked in the conventional servo systems by tracking the point at which both the dissimilar signals are sensed by a servo detector as balanced, in effect relying on the photolithography of the recording elements that generated the servo track edges. An example of a servo track following system for tracking edges is illustrated by the incorporated \"\"430 patent.\nAs pointed out by the incorporated (TU999049) application, tracking at a displacement from an edge is much more difficult. The index position is at a predetermined displacement distance in a lateral direction from an edge. The servo system servo detector determines the ratio between the two dissimilar recorded servo signals read by the servo head, with the ratio, the ratio representing the lateral position of the servo head with respect to an edge. The servo system then moves the servo head laterally to follow the index position, with the servo detector indicating a ratio of the two dissimilar recorded servo signals that is at an offset from the balanced ratio, the ratio representing the desired index servo position. To track follow at a displaced index position requires a high resolution servo detector that can interpolate the dissimilar recorded servo signals. The dissimilar servo signals are prerecorded onto the media, but are subject to variation in amplitude and possibly placement from one data storage media to the next.\nInitializing and calibrating the servo detector to provide a correct interpolation of the recorded dissimilar recorded servo signals therefore becomes difficult.\nAn object of the present invention is to provide an initialization and calibration sequence that enables a more precise track following alignment of a servo head for following an index servo position displaced laterally from an edge.\nDisclosed are a servo system and method for initializing and calibrating at least one index servo track following position substantially parallel to the edges and displaced a predetermined displacement distance in a lateral direction from one of the edges, in accordance with the present invention. A servo head is moveable in the lateral direction with respect to the recording medium and a servo detector is coupled to the servo head for determining a ratio related to the two dissimilar recorded servo signals as read by the servo head, the ratio representing the lateral position of the center of the servo head with respect to one of the edges. A servo track follower is coupled to the servo detector for moving the servo head laterally, the servo track follower, once initialized and calibrated, following an index ratio of the two dissimilar recorded servo signals representing the index servo position. An independent position sensor is provided for determining the mechanical lateral position of the servo head with respect to the recording medium.\nIn accordance with an embodiment of the present invention, logic, coupled to the servo detector, the servo track follower and the independent position sensor, responds to the servo detector and operates the servo track follower to nominally align the servo head at a lateral position at which the servo detector provided ratio represents one of the edges. The logic measures the mechanical lateral position of the independent position sensor at the nominal alignment of the servo head. The logic then responds to the independent position sensor, operating the servo track follower to reposition the servo head laterally the predetermined displacement distance from the nominal alignment as determined by the independent position sensor. The logic measures the servo detector provided repositioned ratio of the servo signals at the displaced distance, the logic initializing and calibrating the servo track follower to employ the provided repositioned ratio as the index ratio. Thus, track following at the provided index ratio insures that the servo head is at the correct displacement from the edge.\nAdditionally, if the servo system comprises a plurality of servo heads having the same alignment as a plurality of the prerecorded servo tracks, the nominal alignment of the servo heads at the corresponding edges comprises aligning the servo heads at a lateral position at which the servo detector provided ratio comprises the average of the provided ratios from the plurality of servo heads, the average representing the corresponding edges of the servo tracks. Then, the measurement of the repositioned ratio of the servo signals comprises measuring the average of the servo detector provided repositioned ratios from the plurality of servo heads.\nIn accordance with another embodiment of the present invention, a check on the edge measurement may be made during the initialization and calibration of index servo positions displaced substantially equidistant in opposite lateral directions from one of the edges. The servo head is again repositioned in the opposite lateral direction from the repositioning step, to the predetermined displacement distance from the nominal alignment as determined by the independent position sensor. The logic measures the again repositioned ratio of the servo signals. Thus, the repositioned ratio represents the displacement at one side of the edge, and the again repositioned ratio represents the displacement at the other side of the edge. The logic determines a midpoint between the repositioned ratio and the again repositioned ratio, compares the midpoint to the nominally aligned ratio, and determines whether the compared ratios are within a predetermined range. If so, the nominal alignment was correct, and the logic initializes and calibrates the servo detector to employ the provided repositioned ratio and the again repositioned ratio as the index ratios for the oppositely displaced index servo positions.\nThe repositioned locations can also be checked in accordance with the present invention. As a check, the servo head is further repositioned laterally in opposite directions from the repositioned displacement, the further repositioning comprising a substantially equal displacement in the opposite directions, the displacement a fraction of the predetermined displacement distance as determined by the independent position sensor. The logic measures the servo detector provided ratios of the servo signals at the further repositioned displacements, and determines a further midpoint between the opposite further repositioned ratios. The logic compares the further midpoint to the repositioned ratio, determining whether the compared ratios are within a predetermined range of each other. If so, the repositioned ratio is correct, and the logic conducts the initializing and calibrating step.\nHence, the present invention provides an initialization and calibration sequence that enables a more precise track following alignment of a servo head for following an index servo position displaced laterally from an edge.\nFor a fuller understanding of the present invention, reference should be made to the following detailed description taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention is related to electronic memories such as may be used for data processors and signal processors.\n2. Description of the Prior Art\nThe prior art uses many different types of memories such as magnetic memories, integrated circuit memories, magnetostrictive delay line memories, and optical memories. Magnetic memories include magnetic core memories and plated wire memories. Integrated circuit memories include shift registers, charge transfer devices (CTDs), and random access memories (RAMs). Bubble memories are magnetic in nature and are implemented with integrated circuit processes. These memories are used with digital systems, where digital information is stored as a single digital bit per memory cell. Hence, a digital word having a plurality of digital bits is stored in a plurality of memory cells storing a single digital bit per memory cell. In analog systems such as analog computers or analog signal processors, analog signals are stored in analog signal form with memories such as potentionmeters and CCD memories. Digital systems do not store digital information in analog signal form, such as interfacing an analog memory to a digital processor with a digital to analog converter at the input to the analog memory and an analog to digital converter at the output from the analog memory interfacing to the processor. Main memories for stored program computers are implemented with digital memories, not analog memories.\nThe prior art is further represented by the art of record in the instant application and in the ancestor applications. Transform processor art and frequency domain processor art is cited in the ancestor applications and is taught in the Rabiner and Gold reference cited hereinafter; in the Goodman, Rabiner and Radner and Bulter and Harvey references cited in ancestor U.S. Pat. Nos. 4,209,853 and 4,209,852; in the Oppenheimer and Schafer reference cited in ancestor U.S. Pat. Nos. 4,209,843 and 4,486,850; and in the Brigham reference cited in ancestor U.S. Pat. No. 4,486,850. Transform processor art and frequency domain processor art is further discussed in the Description Of The Prior Art in ancestor U.S. Pat. No. 4,486,850.\nThe prior art is further defined in the art-of-record of the related applications in the chain of continuing applications including U.S. Pat. No. 3,356,989 to Autry; U.S. Pat. No. 3,613,771 to Quay; U.S. Pat. No. 3,618,052 to Kwei; U.S. Pat. No. 3,643,106 to Berwin; U.S. Pat. No. 3,652,499 to Smith; U.S. Pat. No. 3,662,351 to Ho; U.S. Pat. No. 3,775,738 to Quay; U.S. Pat. No. 3,753,242 to Townsend; U.S. Pat. No. 3,755,793 to Ho; U.S. Pat. No. 3,757,313 to Hines; U.S. Pat. No. 3,761,901 to Aneshansely; U.S. Pat. No. 3,771,148 to Aneshansley; U.S. Pat. No. 3,774,177 to Schaffer; U.S. Pat. No. 3,787,852 to Puckette; U.S. Pat. No. 3,801,967 to Berger; U.S. Pat. No. 3,826,926 to White; U.S. Pat. No. 3,852,745 to Le Bail; U.S. Pat. No. 3,873,958 to Whitehouse; U.S. Pat. No. 3,876,989 to Bankowski; U.S. Pat. No. 3,889,225 to Gosney; U.S. Pat. No. 3,891,977 to Amelio; U.S. Pat. No. 3,895,342 to Mallet; U.S. Pat. No. 3,909,806 to Uchida; U.S. Pat. No. 3,912,748 to Barton; and U.S. Pat. No. 3,942,034 to Buss and including the article by Dennard, IBM Technical Disclosure Bulletin, Vol 14, No. 12, May 1972, pages 3791-3792; the article by Altman, Electronics Magazine, Feb. 28, 1972, pages 62-71; and the article by Baertsch, Electronics Magazine, Dec. 6, 1971, pages 86-91 which references are all incorporated herein by reference."}
{"text": "The present invention pertains to a commissioning device with essentially vertical product storage units, in which products of the same type are stacked, and with at least one conveying means arranged under the product storage units, especially a conveyor belt, for the removal of selected or commissioned products, wherein the product of a selected product storage unit, which product is the lowermost product of the stack, can be pushed out of the product storage unit by an upwardly directed product stop of the conveying means in the direction of conveying of the moving conveying means and can be deposited on a positioned product field of the conveying means.\nA commissioning device of the above-mentioned type has been known from EP 0 960 836 A2. The lowermost product of the stack is removed by the product storage unit being lowered with a continuously driven conveyor belt, which has product stops arranged at spaced locations from one another at the top, to the extent that the product stop or stud moving by carries with it the lowermost product of the column. The drawback of this embodiment is that the weight of the entire product column must be moved to push out a product, which leads to a complicated guide mechanism.\nThe object of the present invention is to provide a commissioning device of the type mentioned in the introduction, which makes possible the reliable and efficient commissioning of at least the lowermost products of stacks of a selected product storage unit by means of simple measures.\nThe essence of the present invention is that the product storage units are arranged lying above the plane of the conveying means, preferably at equally spaced locations from the said plane, and the product stop of the moving conveying means passes through freely under the product storage unit in case of non-commissioning of the lowermost product of the stack and is raised by a lifting means for a lateral engagement with at least one lowermost product of the stack in case of commissioning. Consequently, the entire product storage unit with its full weight is not moved in the direction of the conveyor belt or the product stop to push the product that is the lowermost product of the stack out of the product storage unit, as is done according to the state of the art, but the product stop of the moving conveying means is briefly raised. The product stop can be raised with or without the conveying means. The moving masses are always small, which makes possible a rapid cycling and consequently rapid and reliable commissioning. The moving conveying means may be a conveyor belt, a link conveyor, two toothed belts, V-belts or chains. The advantages of the system according to the present invention are the mechanical simplicity and especially the fact that due to the products being carried with the stop or stud, the products come to lie on the conveying means in an accurately positioned manner, which is favorable, e.g., for automatic labeling. It is also possible to form stacks of products on the conveying means between two adjacent product stops located at spaced locations from one another without an additional effort.\nProvisions are made, in particular, for the conveying means to be raised by the lifting means together with the product stop rigidly connected to the conveying means directly in front of and after the product of a product storage unit, which is the lowermost product of the stack, and to be lowered again into the plane of the conveying means after pushing out.\nThe lifting means are preferably a cycled, vertical lifting drive, which is provided under the conveying means in the area of an associated product storage unit and is operated especially pneumatically or electrically.\nThe lifting means may also include an oblique stationary ramp with a wedge-shaped lifting body and be provided under the conveying means in the area of an associated product storage unit, the lifting body, which is displaceable or can be moved by means of rollers, being able to be adjusted in height along the ramp.\nAs an alternative, the conveying means may be, in principle, longitudinally movable, but not adjustable in height, and the pushing out of a product that is the lowermost product of the stack may also be pushed out of the selected product storage unit by the product stop connected to the conveying means being raised by lifting means without the conveying means directly in front of and after the product of a product storage unit, which product is the lowermost product of the stack and is to be commissioned, and being again lowered into the plane of the conveying means after pushing it out.\nThe product stop is preferably connected in this case to the conveying means by means of a conveying means tab, which is subjected to tensile load, and the product stop is adjustable in height by means of literal stationary ramps in the area of the lateral edge of the conveying means or by a means arranged above the conveying means at the product storage unit when commissioning is to be performed.\nThe lateral ramps advantageously have an adjustable, cycled switch each, which is preferably a laterally displaceable part of the ramp. If the switch is inactive, a transverse connecting rod of the product stop can be moved in the direction of conveying in a lateral longitudinal groove of the lateral ramps at the same level as the conveying means. If the switch is set, the longitudinal groove is blocked by slides, the transverse connecting rod passes through a predetermined elevated path corresponding to the preset ascent of the ramp, the saddle of the ramp and the decline of the ramp and consequently also corresponding to the product stop, which will then laterally engage a lowermost product of the stack and pushes out this product in the course of its forward movement and positions same on the product field in front of the product stop on the conveying means.\nInstead of a preferably central product stop, two lateral product stops of identical design may also be provided for a single product of a product storage unit, which is the lowermost product of the stack, in order to reliably prevent jamming from occurring during the pushing out in all cases.\nThe product storage units may be preferably slightly sloped relative to the vertical in the direction of conveying of the conveying means. The pushing-out operation is facilitated by the slope.\nThe product storage units of a conveying means may be arranged in a row and spaced equally from one another.\nIn particular, the distance between the product storage units is smaller than the width of the product storage unit in order to ensure the compact design of the overall system. The distance between two product storage units can be bridged over by a protective plate in order to minimize the risk for injury to a user. A high product stop moving forward may possibly represent a safety hazard.\nThe product storage units may also be sloped, preferably slightly, relative to the vertical at right angles to the direction of conveying of the conveying means.\nProvisions may be made, in particular, for the product storage unit to be a vertical, doubly sloped angle sheet iron with a positioning angle for products picked up in the product storage unit. Differently configured cubic objects can also be stacked in an aligned manner in such a storage unit and handled and especially ejected correspondingly reliably.\nThe product storage units may have a lateral ejection slot located in the direction of conveying for products to be given out, where the width of the slot is selected or the slot can be adjusted in width in the vertical direction such that two or more stacked products can be deposited simultaneously as a small stack of products on a single product field of the conveying means. It is obvious that the height of the product stop must be selected correspondingly in this case. A product stop may be optionally moved in height on two different lateral ramps or on a single lateral ramp, which ramp will now have two switches of different heights.\nIn particular, an ejection slot that is broad in the vertical direction may be formed in the product storage unit, which has a maximum vertical extension for pushing out, e.g., two or three products at the same time, with a brush being fastened horizontally in the upper horizontal edge area of the ejection slot such that the bristles of the brush extend downward and end under the top edge of the product of the product storage unit, which product is the lowermost product in the stack, and a product falling after is reliably retained. If only the product that is the lowermost product of the stack is commissioned in the case of such a design, this product can be pushed out practically freely with the product stop touching the bristles only slightly. If two (or three) products are being commissioned, and the raised product stop has a correspondingly high engagement position, the second lowest (or possibly also the third lowest) product of the product stack is also pushed out of the product storage unit together with the lowermost product of the stack. The elastic bristles of the brush are correspondingly deflected in this case by the second lowest (and possibly the third lowest) product during pushing out. It is obvious that another elastic member, e.g., leaf spring tongues, which performs the same function, may also be provided instead of the brush.\nInstead of the brush or the elastic member, it is also possible to provide a plate slide, which releases the entire vertical width of the ejection slot in a raised position, so that, e.g., two or three products can be pushed out simultaneously. If only the lowermost product of the stack is to be pushed out or commissioned, the height-adjustable slide is correspondingly lowered in order to release only the width [necessary] for the lowermost product of the stack. The plate slide is correspondingly actuated now by an electronic control means, which is part of the electronic control unit of the commissioning device, just as the product stop for its height adjustment.\nThe overall system of a commissioning device may have different designs depending on the capacity, the intended use and the local conditions. In particular, a plurality of conveying means with associated product storage units may be provided in the same plane and/or in different planes, whose conveying end is in connection with a cross belt or the like, and oblique intermediate conveyor belts or the like, which can be bridged over in height, may be provided.\nA plurality of pushed-out products, which have been taken from one or more product storage units, may be preferably stacked in a single product field.\nOne or more product storage units may be designed as drawers or as roller-mounted drawers displaceable at right angles to the conveyor belt. The product storage unit is now pulled out laterally for filling. The position of the drawer is scanned by means of a sensor or scanner. Pushing out is prevented by the product storage unit control for the time of refilling. If a product is needed from the blocked product storage unit during this time, the product storage unit control reports the error or the refilling operation to the control computer.\nThe drawer may have a plurality of product storage units next to one another (parallel channels). The unused product storage units can be filled during the operation. A full product storage unit is then brought into the push-out position by displacing the drawer, and the empty product storage unit can again be refilled. The pushing out must be prevented in this case only for the duration of the channel change. The product storage unit can be correspondingly enlarged at equal overall height by the automatic or manual pushing on of a drawer with a plurality of parallel, channels.\nConsequently, a stack of products (product column) is present according to the present invention in the product storage unit. The product storage units are always arranged vertically or sloped (at right angles to the direction of conveying) above a lifting device located under the conveyor belt).\nThe studs or product stops have such a height that they can pass through under the magazines with the lifting device withdrawn and the belt moving. If a product is to be removed, the conveyor belt is raised with the belt moving at the moment at which a stud is located directly in front of the front edge of the product storage unit. The conveyor belt is raised to the extent that the stud carries the lowermost product of the product column in the direction of conveying and as a result pushes it out of the product storage unit. The pushed-out product will then be located on the conveyor belt in a discrete field, whose size is determined by the distance between the studs.\nA plurality of products can be pushed out of different product storage units with the same studs. As a result, a product stack is formed, which can be taken over at the end of the conveyor belt.\nThe advantages of the system are the mechanical simplicity and the fact that due to the products being carried with the studs, the products lie on the belt in an accurately positioned manner (which is favorable, e.g., for automatic labeling), and product stacks can be formed without an additional effort.\nEspecially in the area of E-commerce, where an order consists, in general, of a small number of products only, a complete order requires only one product window if the height of the stack does not exceed the height of the stud. A drastic increase in the throughput of the system is achieved as a result.\nThe raising of the conveyor belt, or the like can be accomplished in many ways. One possibility is to push a plastic wedge over a ramp recessed in the material handling equipment. Any method for generating a lifting movement is possible, in principle (lifting magnet, vertically installed cylinder with roller, lateral cylinder with lifting mechanism, etc.).\nControl:\nThe control computer receives an order consisting of a plurality of products. If the height of the stack is smaller than the height of the stud, the entire order can be stacked into a discrete field (product window=distance between two studs). If an order requires a plurality of product windows, the order is split. Due to the arrangement of the product storage units with lifting mechanism above the materials handling equipment, each storage unit has a defined three-dimensional distance from the end (in the direction of conveying) of the materials handling equipment. A pulse generator coupled with the materials handling equipment provides the information on the position of the product window. The product windows needed for an order are first reserved. If the product window is then located under a product storage unit, from which a product is to be pushed out, the corresponding lifting mechanism is activated.\nTo save space, a plurality of planes are connected one over another and/or by a curve and/or are arranged one behind the other, but also next to one another. To again reach a product window after a curve, the velocity of the conveyor belts of this plane must be equal, and the phase position of the product windows of the plane must remain constant to one another, i.e., it must be able to be configured. Finally, all the planes located one on top of another and one behind the other and optionally also next to one another are brought together onto one section. To reach the product windows of the collecting belt during this operation, the velocity and the phase position of the converging sections must also be constant or able to be configured. The belt synchronization control is necessary only in the case of curves and to optimize the throughput.\nAs an alternative, the raising of the stud may also take place according to another principle. Instead of raising the conveyor belt, a conveyor belt or the like with a tab, strip, small plate or the like, arranged on the conveying side, may be used. The tabs or the like are fastened at the conveyor belt at right angles to the direction of conveying and have a stud at the loose end, which stud acts as a push-out means. Together with the stud, a rod, which projects over the lateral edge of the conveyor belt, is mounted at right angles to the direction of conveying. A ramp with a switch is located under each product storage unit on both sides next to the conveyor belt. If a product is to be pushed out, the switch is activated as soon as the tab belonging to the product window is located under the product storage unit. The product storage units are sloped in the direction of conveying. In addition, the product storage units may also be sloped at right angles to the direction of conveying in order to make possible the simple filling of the product storage unit, because the products are automatically in contact with the rear side due to the force of gravity, and the front side (seen at right angles to the direction of conveying) can be left open. The rod at the end of the tab runs up onto the ramp and is raised as a result. After overcoming the ascent, there is a gradient, which corresponds to the slope of the product storage unit. The knob raised by the rod now carries the lowermost product with it out of the product storage unit. The slope of the product storage unit is selected to be such that the pushed-out product can be pushed under the next channel. As a result, the distance between consecutive product storage units can be minimized. If the switch is not activated, the rod passes through under the ramp along a groove.\nA combination of the two possibilities is also conceivable: A higher stud with a rod at right angles to the direction of conveying is raised over a ramp only to the extent that the top edge of the stud pushes the lowermost product out of the product storage unit. The possibility of forming stacks is thus preserved. However, the ramp may also be part of the product storage unit, and projections on the stud are used to raise the stud.\nOverall system:\nThe system may comprise a plurality of planes located one on top of another and rows located next to one another. A container handling equipment is arranged at right angles to the rows. The products are transferred at the end of the materials handling equipment either directly in each plane into containers waiting in front of it, or the planes are converged onto one level with transfer of the products into containers in one level.\nThe various features of novelty which characterize the invention are pointed out with particularity in the claims annexed to and forming a part of this disclosure. For a better understanding of the invention, its operating advantages and specific objects attained by its uses, reference is made to the accompanying drawings and descriptive matter in which preferred embodiments of the invention are illustrated."}
{"text": "Unless otherwise indicated herein, the materials described in this section are not prior art to the claims in this disclosure and are not admitted to be prior art by inclusion in this section.\nA news production system (NPS) may facilitate the production of a news program in the form of a media stream. In one example, an NPS may include multiple media sources and a production switcher, where outputs of the media sources are connected to inputs of the production switcher. This may allow the production switcher to switch between and/or combine multiple media streams output by the media sources, thereby outputting the news program in the form of another media stream.\nThere are various types of media, including for example, audio, video, or a combination thereof. As such, in one example, an NPS may output a news program in the form of an audio stream. In this instance, the NPS may transmit the audio stream to a radio-broadcasting system for broadcast. As another example, a media stream may take the form of a video stream or a combined audio and video stream. In such instances, the NPS may transmit the video stream or the combined audio and video stream to a television-broadcasting system for broadcast.\nA media source may take a variety of forms. For example, a media source may take the form of a media server. A media server is a device configured for retrieving a media file, converting the retrieved media file into a media stream, and outputting the converted media stream.\nAs another example, a media source may take the form of a media effect engine. A media effect engine is a device configured for retrieving a media effect (sometimes referred to as a “page”), and running the media effect thereby outputting a corresponding media stream. A media effect may be stored as a file that includes instructions and other data (e.g., media) related to the media effect. By running the media effect, the media effect engine may generate and output a media stream based on those instructions. Media effects are commonly used as a means to generate animations, graphics, or other visual effects in the form of a media stream that can be overlaid on another media stream. For instance, in the context of a news program, a “lower third” media effect may be used to overlay a graphic over a lower third portion of a media stream.\nAs such, in one example NPS, a media server may output a first media stream while a media effect engine outputs a second media stream, and a production switcher may combine the two media streams (e.g., by overlaying the second media stream over the first media stream) to output the news program in the form of a third media stream. The production switcher may then transmit the third media stream to a broadcasting system (e.g., a television-broadcasting system) for broadcast.\nA media effect engine may be controlled in a variety of ways such that it may perform the steps of retrieving a media effect and running the media effect. For instance, a user may control a media effect engine by providing it a suitable instruction via a user-interface. However, for a variety of reasons, this manner of controlling the media effect engine may be undesirable. Among other things, this process may be time-consuming for the user. In addition, it may be difficult for the user to ensure that the media effect engine performs such steps at appropriate times during production of the news program.\nAs another example, a controller device may control a media effect engine by providing it a suitable instruction in accordance with one or more application programming interfaces (API) that may be made publically available by the provider of the media effect engine or another entity. However, again for a variety of reasons, this manner of controlling the media effect engine may be undesirable. Among other things, it may be time-consuming for a user to configure the controller device to provide such a suitable instruction.\nThis approach may be particularly time-consuming given that different instructions may need to adhere to different APIs. For instance, an instruction requesting the running of one type of media effect may need to adhere to a different API than another instruction requesting the running of another type of media effect. In addition, it may be difficult for the user to configure the controller device such that it causes the media effect engine to retrieve and run a particular media effect at an appropriate time during production of the news program."}
{"text": "A. Field of Invention\nThe present invention relates generally to methods of constructing and controlling moving telescoping boom sections of an aerial lift or other device and a telescoping boom constructed in accord with the method, and more particularly to a new and improved method of constructing and controlling moving telescoping boom sections of an aerial lift or other device, each of the boom sections being extended or retracted by a hydraulic cylinder and a new and improved telescoping boom apparatus constructed in accord with said method, having, for each moving section, a rigidly affixed hydraulic cylinder.\nB. Description of Related Art\nMany types of aerial lifts, cranes and similar telescoping boom devices have been provided. Further, many of these use a hydraulic actuator to extend or retract the boom. However, it is not believed that any of these devices use an individual hydraulic actuator for each moving section, such lifts instead conventionally using chains, cables, or the like in pulley systems to achieve the extension and retraction of the boom sections. Nor do conventional lifts have rigidly mounted hydraulic actuators or sliding supports on the hydraulic actuator piston rods."}
{"text": "A transparent electrode is generally used for: a liquid crystal display, an electroluminescent display, a plasma display, an electrochromic display, a solar battery, a touch panel, and an electronic paper.\nAn organic EL element (it may be called as an organic-field light-emitting element), which utilizes electroluminescence (hereinafter, it is abbreviated as “EL”) of an organic material, has a configuration of interposing a light emitting layer containing an organic compound between a pair of opposed electrodes. Emission light generated in the light emitting layer passes through the electrode and it is extracted to the outside. Therefore, at least one of the two electrodes is composed of a transparent electrode.\nAs for a transparent electrode, oxide semiconductor materials, such as indium tin oxide (SnO2—In2O3: or abbreviation name ITO), are generally used. The transparent electrode made of ITO is usually produced with a sputtering method. The transparent electrode produced with a sputtering method only will have a large sheet resistance, and it exhibits remarkable voltage decrease from the power supplying point.\nIn order to decrease a sheet resistance, it was investigated a method of laminating an ITO layer and a silver layer for reducing resistance (for example, refer to Patent documents 1 and 2).\nPatent document 1 discloses a transparent electrode having a structure of laminating an ITO film and a silver film. Patent document 2 discloses a transparent electrode having a structure of interposing a silver film with ITO films. However, when the used silver film was made thin to an extent of not inducing loss of light transmittance, the resistance did not sufficiently decrease. Therefore, it was required to combine with a metal oxide such as ITO. Since ITO employs a rare earth metal In, the cost of material is high. In addition, it is required to perform an annealing treatment at a temperature condition of around 300° C. for decreasing the resistance. It was difficult to use a resin substrate at such temperature condition. When silver was employed, high electric conductivity is obtained. On the other hand, it has a problem of trade-off between resistance and light transmittance.\nPatent document 3 proposes the structures for replacing a metal oxide such as ITO. One structure contains a thin film metal material such as silver having high electric conductivity. The other structure contains a mixture of silver with aluminum. This mixture enabled to achieve higher electric conductivity with a thinner film than silver. However, a transparent electrode composed of silver and aluminum having high electric conductivity had slightly insufficient sheet resistance. A metal thin film is usually is required to have a large thickness to increase sheet resistance. Consequently, light transmittance will be deteriorated. As described above, the resistance property and the light transmittance are in a trade-off relationship. It was difficult to achieve sufficient conductivity and light transmittance at the same time. Further, in a silver-aluminum alloy, aluminum has a property of easily oxidized. It has a problem of increasing the resistance by oxidization. Therefore, it has been required a transparent electrode enabling to achieve both resistance (conductivity) and light transmittance."}
{"text": "1. Field of the Invention\nThe invention relates to a control system for an automatic transmission.\n2. Description of the Related Art\nIn a conventional automatic transmission, rotation from the engine is transmitted to a speed change unit through a hydraulic power transmission, and speed changes are effected in the speed change unit. A first clutch (an input clutch) is arranged between the hydraulic transmission and the speed change unit for change between a neutral range and a forward driving range by engaging/disengaging the first clutch.\nIn the automatic transmission, a neutral control is performed by disengaging the first clutch based on a hydraulic pressure modulated by a linear solenoid valve when the forward driving range is selected, an accelerator pedal is released, a brake pedal is depressed and the vehicle is in a \"vehicle stopped state\" defined as a vehicle speed of almost 0. Fuel consumption is improved by such neutral control because the load on the engine is reduced and vibration of the engine is prevented.\nWithin the neutral control, a hill-hold function is provided so that the vehicle will not move backward when the brake pedal is released for starting, after the vehicle has been stopped on an incline and facing up-hill. To provide the hill-hold function, when it is detected that the vehicle is in the \"vehicle stopped state\", a hill-hold output is output to establish a hill-hold state in the speed change unit by, for example, engaging a first brake and locking a designated one-way clutch.\nHowever, in the conventional automatic transmission, when the first brake is engaged very shortly after the hill-hold output, the hill-hold state is established while engaging the first clutch. In this case, the automatic transmission is rotated to a small degree, in one direction around an input shaft, by a reaction force received before the hill-hold state is established and held. When the first brake is engaged before disengaging the first clutch in this state, the automatic transmission is held in a state wherein the automatic transmission receives the reaction force. The engine, connected to the automatic transmission, is thereby held in the same state. As a result, the neutral control is performed in a state wherein the engine is under a load. Therefore, the load on the engine increases accompanying the change in the conditions of the neutral control. As a result, the fuel consumption is not improved and vibration of the engine is not reduced to the extent that it cannot be felt in the driver's seat."}
{"text": "The present invention relates to a device and a method for dispensing a substance, for example, in one use or application, for self-administering a product fluid, e.g., insulin.\nThe invention also relates to a container for a substance, in particular a product fluid to be self-administered, and to a device for administering a substance.\nDevices for dispensing substances wherein the device is adapted and used for a specific application or use, such as insulin treatment, at a particular dose, are known. In such cases, due to the specific application modalities, it is necessary to produce different dispensing devices for different areas of application having different task specifications.\nSome devices for administering product take the form of portable infusion and injection devices, some of which are used in insulin treatment. Generally, such devices, including those used in insulin treatment, involve containers filled with the substance to be administered which are coupled to an administering device in order to dispense the substance contained in the container via the administering device to a patient. There are a multitude of substances which can be administered in this way, such as preparations comprising insulin for diabetes, growth hormones (hGH; human Growth Hormone) for disturbed growth, erythropoietine (Epo) for renal insufficiency or general lack of red blood corpuscles, alpha-interferone for hepatitis or cancer treatment, or potency-stimulating agents. Such containers, which are often geometrically identical, are often filled with different concentrations of the substance to be administered.\nIn order to reduce the danger of confusing containers having different substances, variously formed administering devices are known into which the respectively corresponding containers can be inserted.\nA container is known from WO 98/00187 comprising a color coding which can be attached to it, consisting of a number of variously colored fields, wherein a property of a container or its contents can be identified by means of an optical sensor system."}
{"text": "The present invention relates to exercise equipment, and more particularly to exercise treadmills, and still more particularly to suspension systems for supporting the deck of the exercise treadmill above an underlying frame structure.\nExercise treadmills are widely used in spas, exercise clubs and also in individual residences to enable users to walk, jog or run indoors. This is especially useful during inclement weather and also at night or at other times when exercisers do not desire to run outdoors. Most exercise treadmills include first and second roller assemblies that are transversely mounted at the ends of a frame. An endless belt is trained about the roller assemblies. The upper run of the belt is supported by an underlying deck positioned between the belt and the frame.\nEfforts have been made to reduce the impact on the user\"\"s limbs and joints when jogging or running on a treadmill. One method of reducing the impact on an exerciser\"\"s body is disclosed by U.S. Pat. Nos. 4,974,831 and 4,984,810. In the treadmills disclosed by these patents, the rear end of the deck is pivotally mounted to the frame, with the forward end of the deck supported by a suspension system. In the \"\"831 patent, the suspension system consists of a fairly complicated lever arm assembly and cooperating shock absorbers. Striding on a deck results in pivoting of the lever arms and extension of the shock absorbers, thereby to dampen the impact of the user\"\"s feet. A drawback of this shock absorption system is its complex nature, rendering it costly to manufacture.\nIn the \"\"810 patent, the forward end of the treadmill deck was supported by a conventional compression spring and separate shock absorber. Placement of the spring and shock absorber at the very front of the deck imposes considerable bending stress on the deck.\nOther conventional treadmills have utilized rubber blocks positioned between the deck and the underlying frame to absorb impact. One such conventional treadmill is disclosed in French Patent No. 2,616,132. A treadmill deck is mounted above the frame members on a plurality of flexible pads. Bushings are inserted into the top and bottom of each pad, and bolts depending downwardly from the deck and upwardly from frame are received within the corresponding bushings. The bolts serve to position the flexible pads between the deck and frame for shock absorption.\nU.S. Pat. Nos. 5,336,144 and 5,454,772 disclose a deck supported above a frame by a plurality of cup-shaped elastomeric springs. The elastomeric springs reversibly deform during downward deflection of the deck toward the frame. The elastomeric springs have side walls of tapering thickness. As a result, the resistance to the downward travel of the deck provided by the elastomeric springs is proportional to the degree of deflection of the deck toward the frame. One drawback of this particular treadmill construction is that the elastomeric springs are fixed in place and individually define a rather small bearing area.\nThe present invention provides an exercise treadmill having a frame, first and second roller assemblies rotatably mounted on the frame, and an endless belt trained about the first and second roller assemblies. The exercise treadmill also includes a deck disposed between the frame and the upper run of the belt. A pivot connection pivotally connects the rearward end portion of the deck to the frame. Elongate elastomeric spring members are disposed between the frame and the deck at a location intermediate the ends of the deck to support the deck spaced above the frame. The elastomeric springs reversibly deform to resist a deflection (downward movement) of the deck toward the frame when the exerciser strides on the endless belt. The resistance provided by the elastomeric spring members is proportional to the extent of deflection of the deck.\nIn a further aspect of the present invention, the elastomeric spring members are mounted on the side rails of the frame and underlie marginal side portions of the deck.\nIn another aspect of the present invention, the elastomeric springs include a base portion and a bulbous body portion extending upwardly from the base portion. The body portion is domed or crowned at its top to define an outwardly convex shape. The interior of the elastomeric spring between the base portion and the body portion is hollow or partially hollow. As a result, the body portion deflects downwardly under the force imposed on the deck by the exerciser.\nIn an additional aspect of the present invention, the wall thickness of the body portion of the elastomeric spring is greater at the intersection of the body portion with the base of the elastomeric spring. The wall thickness of the body portion decreases in the direction away from the base portion, reaching a minimum thickness at the top of the domed body portion. As a result, when the deck imparts a downward load on the elastomeric springs, the top central portion of the body portion of the elastomeric spring deflects downwardly into the hollow interior, rather than the body portion deflecting sideways, which could occur if the elastomeric spring was of solid construction. Also, the resistance imposed on the deck by the elastomeric spring increases as the deck deflects downwardly, thereby providing a variable rate spring.\nIn another aspect of the present invention, the spring may be constructed so that its rate of deformation may be selectively altered. In this regard, a compressible insert is sized and shaped to be selectively insertable to a desired degree into the hollow body portion of the spring. In cross-section, the insert may correspond to the cross-sectional shape of the hollow body portion of the spring. Also, the spring may be tapered along its length. In another configuration, the body portion of the spring may be adapted to receive a compressible fluid thereby serving as a bladder. In a more specific aspect of the present invention, the compressible fluid may be composed of air, with the air being supplied to the bladder by an air pump. Also in a more specific aspect of the present invention, a valve or other means may be provided for discharging the compressible fluid from the bladder.\nIn a further aspect of the present invention, the pivot connection at the rearward end of the deck includes a spindle mounted on the frame side member to engage with a hinge bracket mounted to the underside of the deck. By this construction, the rearward end portion of the deck is pivotally attached to the frame about an axis extending transversely to the length of the deck."}
{"text": "1. Field of the Invention\nThe present invention relates to a predistortion compensation apparatus for performing distortion compensation processing in advance to a transmission signal before amplification.\n2. Description of the Related Art\nIn recent years, high-efficient digital transmission has been adopted in the radio communication field. When multilevel phase modulation is adopted in the radio communication, a technique for reducing adjacent channel leak power becomes important, in which nonlinear distortion is restrained by linearizing the amplification characteristic of a power amplifier on the transmission side.\nAlso, to improve power efficiency even in case an amplifier having a degraded linearity is used, a technique for compensating nonlinear distortion for the degraded linearity is necessary.\nFIG. 1 shows an exemplary block diagram of transmission equipment in the conventional radio equipment. A transmission signal generator 1 outputs a digital serial data sequence. Also, a serial-to-parallel (S/P) converter 2 converts the digital data sequence into two series, in-phase component (I-component) signals and quadrature component (Q-component) signals, by alternately distributing the digital data sequence on a bit-by-bit basis.\nA digital-to-analog (D/A) converter 3 converts the respective I-signals and Q-signals into analog baseband signals, and inputs the signals into a quadrature modulator 4. This quadrature modulator 4 performs orthogonal transformation and outputs signals by multiplying the input I-signals and Q-signals (transmission baseband signals) by a reference carrier wave 8 and a carrier wave phase-shifted therefrom by 90°, respectively, and adding the multiplied results.\nA frequency converter 5 mixes the quadrature modulation signals with local oscillation signals, and converts the mixed signals into radio frequency. A transmission power amplifier 6 performs power amplification of the radio frequency signals output from frequency converter 5, and radiates the signal to the air from an antenna 7.\nHere, in the mobile communication using W-CDMA, etc., transmission equipment power is substantially large, becoming as much as 10 mW to several tens of mW, and transmission power amplifier 6 has a nonlinear input/output characteristic having a distortion function f(p), as shown by the dotted line in FIG. 2. This non-linearity causes a non-linear distortion. As shown by the solid line (b) in FIG. 3, the frequency spectrum in the vicinity of a transmission frequency f0 comes to have a raised sidelobe from the characteristic shown by the broken line (a). This leaks to adjacent channels and produces adjacent interference. Namely, due to the nonlinear distortion shown in FIG. 2, leak power of the transmission wave to the adjacent frequency channels becomes large, as shown in FIG. 3.\nAn ACPR (adjacent channel power ratio) is used to indicate the magnitude of leak power. ACPR is a ratio of leak power to adjacent channels to the power in the channel of interest, in other words, a ratio of the spectrum area in the adjacent channels sandwiched between the lines B and B′ in FIG. 3 to the spectrum area between the lines A and A′. Such leak power affects other channels as noise, and degrades communication quality of the channels concerned. Therefore, a strict regulation has been established to the issue of leak power.\nThe leak power is substantially small in a linear region of, for example, a power amplifier (refer to a linear region I in FIG. 2), but is large in a nonlinear region II. Accordingly, to obtain a high-output transmission power amplifier, the linear region I has to be widened. However, for this purpose, it becomes necessary to provide an amplifier having a larger capacity than is actually needed, which causes disadvantage in apparatus cost and size. As a measure to solve this problem, a distortion compensation function to compensate for transmission power distortion is added to radio equipment.\nFIG. 4 shows the block diagram of transmission equipment having a digital nonlinear distortion compensation function by use of a DSP (digital signal processor). A digital data group (transmission signals) transmitted from transmission signal generator 1 is converted into two series, I-signals and Q-signals, in S/P converter 2, and then the two series of signals are input to a distortion compensator 9.\nAs shown in the lower part of FIG. 4 in enlargement, distortion compensator 9 includes a distortion compensation coefficient storage 90 for storing a distortion compensation coefficient h(pi) corresponding to the power level pi (i=0-1023) of a transmission signal x(t); a predistortion portion 91 for performing a distortion compensation process (predistortion) onto the transmission signal, using the distortion compensation coefficient h(pi) corresponding to the transmission signal power level; and a distortion compensation coefficient calculator 92 for comparing the transmission signal x(t) with a demodulation signal (a feedback signal) y(t) demodulated in the quadrature detector which will be described later, and calculates and updates the distortion compensation coefficient h(pi) so that the difference between the transmission signal and the demodulation signal becomes zero.\nThe signal to which distortion process is performed in distortion compensator 9 is input into D/A converter 3. D/A converter 3 converts the input I-signal and Q-signal into analog baseband signals, and inputs the converted signals into quadrature modulator 4. Quadrature modulator 4 performs quadrature modulation by multiplying the input I-signal and Q-signal by a reference carrier wave 8 and a carrier wave being phase-shifted from carrier wave 8 by 90°. Quadrature modulator 4 then adds and outputs the multiplied result.\nA frequency converter 5 mixes the quadrature modulation signal with a local oscillation signal, and performs frequency conversion. A transmission power amplifier 6 performs power amplification of the radio frequency signal output from frequency converter 5, and radiates the signal to the air by an antenna 7.\nA portion of the transmission signal is input to a frequency converter 11 via a directional coupler 10, and input into a quadrature detector 12 after being converted by the above frequency converter 11. Quadrature detector 12 performs quadrature detection by multiplying the input signal by a reference carrier wave, and by a signal which is phase shifted by 90° from the reference signal, respectively. Thus, the baseband I-signal and Q-signal on the transmission side are reproduced, which are then input into an analog-to-digital (A/D) converter 13.\nA/D converter 13 converts the input I-signal and Q-signal into digital signals, and inputs into distortion compensator 9. Through the adaptive signal processing, using an LMS (least-mean-square) algorithm, in distortion compensation coefficient calculator 92 of distortion compensator 9, the pre-compensated transmission signal is compared with the feedback signal being demodulated in quadrature detector 12. Then distortion compensator 9 calculates the distortion compensation coefficient h(p1) so as to make the above difference zero. Then, distortion compensator 9 updates the above-obtained coefficient which has been stored in distortion compensation coefficient storage 90. Through the repetition of calculations above, nonlinear distortion in transmission power amplifier 6 is restrained, and adjacent channel leak power is reduced.\nBy way of example, in the PCT International Publication WO 2003/103163, such a configuration as shown in FIG. 5, in which distortion compensation is performed using the adaptive LMS algorithm, is described as an embodiment of distortion compensator 9 shown in FIG. 4.\nIn FIG. 5, a multiplier 15a corresponds to a predistortion section 91 shown in FIG. 4, in which a transmission signal x(t) is multiplied by a distortion compensation coefficient hn-1(p). Also, a distortion device 15b having a distortion function f(p) corresponds to a transmission power amplifier 6 shown in FIG. 4.\nFurther, as to the portion in FIG. 4 including a frequency converter 11, a quadrature detector 12 and an A/D converter 13, in which the output signal being output from transmission power amplifier 15b is feedbacked, a feedback system 15c is shown in FIG. 5.\nMoreover, in FIG. 5, a look-up table (LUT) 15e constitutes a distortion compensation coefficient storage 90 shown in FIG. 4. A distortion compensation coefficient calculation section 16 constitutes a distortion compensation coefficient calculation section 92 shown in FIG. 4, which generates an update value of the distortion compensation coefficient stored in look-up table 15e. \nIn the distortion compensation apparatus having the configuration shown in FIG. 5, look-up table 15e has a distortion compensation coefficient for canceling the distortion produced in transmission power amplifier 6, namely, distortion device 15b, in a two-dimensional address location corresponding to each discrete power value of the transmission signal x(t).\nWhen the transmission signal x(t) is input, an address generation circuit 15d calculates the power p (=x2(t)) of the transmission signal x(t), and generates an address of one dimensional direction, for example the X-axis direction, which uniquely corresponds to the above-calculated power p (=x2(t)) of the transmission signal x(t). At the same time, address generation circuit 15d obtains a difference ΔP of the power P1 (=x2(t-1)) of the transmission signal x(t-1) of the previous time point (t-1) having been stored in address generation circuit 15d, and generates an address of the other dimensional direction, for example, the Y-axis direction, which uniquely corresponds to the above difference ΔP.\nThus, from address generation circuit 15d, a store location in look-up table 15e, which is specified by the address P in the X-axis direction and the address ΔP in the Y-axis direction, is read out. The readout address is output as address designation information (AR).\nThen, a distortion compensation coefficient hn-1(p) stored in the above readout address is read out from look-up table 15e, so as to be used in the distortion compensation processing performed by multiplier 15a. \nMeanwhile, an update value for updating the distortion compensation coefficient having been stored in look-up table 15e is calculated in a distortion compensation coefficient calculation section 16. Namely, distortion compensation coefficient calculation section 16 is constituted of a conjugate complex calculation section 16 and multipliers 15h-15j. A subtractor 15g outputs a difference e(t) between the transmission signal x(t) and the feedback demodulation signal y(t). Multiplier 15i multiplies the distortion compensation coefficient hn-1(p) by y*(t), so as to obtain an output u*(t) (=n-1(p)y*(t)). Multiplier 15h multiplies the difference e(t) being output from subtractor 15g by u*(t). Multiplier 15j multiplies the output of multiplier 15h by a step size parameter μ.\nNext, an adder 15k adds the distortion compensation coefficient hn-1(p) to the output μe(t)u*(t) being output from multiplier 15j, and obtains an update value of look-up table 15e. This update value is to be stored in a write address (AW), consisting of the X-axis direction address and the Y-axis direction address, being specified by address generation circuit 15d as the address corresponding to the transmission signal power p (=x2(t)).\nHere, the aforementioned write address (AW) is the same address as the readout address (AR). However, because of a calculation time, etc. needed to obtain the update value, the readout address is used as the write address after the readout address is delayed in a delay section 15m. \nDelay portions 15m, 15n, 15p add to the transmission signal x(t), the delay time D, which is the period from the input of the transmission signal x(t) to the feed back decoded signal y(t) input to the subtractor 15g. \nThe delay time D being set by the delay portions 15m, 15n, 15p is determined so as to satisfy D=D0+D1, where D0 is the delay time in transmission power amplifier 15b, and D1 is the delay time in feedback system 15c. \nUsing the above configuration, the following calculations are performed.hn(p)=hn-1(p)+μe(t)u*(t)e(t)=x(t)−y(t)y(t)=hn-1(p)×(t)f(p)u*(t)=x(t)f(p)=hn-1(p)y*(t)p=|x(t)|2 Here, x, y, f, h, u, e are complex numbers, and * denotes a conjugate complex number.\nThrough the above calculation processing, the distortion compensation coefficient h(p) is updated so as to minimize the differential signal e(t) between the transmission signal x(t) and the feedbacked demodulation signal y(t). Finally, the value converges to an optimal distortion compensation coefficient, so that the distortion of the transmission power amplifier is compensated.\nNow, in the above calculation, the step size parameter μ determines a degree of effect of an error component e(t) between the transmission signal x(t), i.e. the reference signal, and the feedback demodulation signal y(t), i.e. the feedback signal, on the update value of the distortion compensation coefficient. In the conventional system, the value of the step size parameter μ is set to a fixed value.\nIn the configuration of the distortion compensation apparatus shown in FIG. 5, the inventor of the present invention has observed the output of the distortion compensation apparatus by inputting the outputs of transmission power amplifier 15b, the distortion device, into a spectrum analyzer with sweep frequencies. FIGS. 6A through 7B are the results obtained at those times. In the examples shown in FIGS. 6A through 7B, the observations have been made with different transmission signal levels in four frequency bands (channels).\nFIGS. 6A, 6B represent the output spectrum waveforms of the spectrum analyzer when the transmission signal level is large (43 dB). FIG. 6A shows the spectrum waveform when the step size parameter μ is set to 1/1024, while FIG. 6B shows the spectrum waveform when the step size parameter μ is set to 1/16.\nAlso, FIGS. 7A, 7B represent the output spectrum waveforms of the spectrum analyzer when the transmission signal level is small (27 dB). FIG. 7A shows the spectrum waveform when the step size parameter μ is set to 1/1024, while FIG. 7B shows the spectrum waveform when the step size parameter μ is set to 1/16.\nFrom these FIGS. 6A through 7B, it has been found out that the relation between the transmission signal level and the step size parameter μ produces an effect on the distortion compensation coefficient. In FIGS. 6A, 6B, if the step size parameter μ is set larger when the transmission signal level is large (refer to FIG. 6B), external disturbance (phase rotation, quantization error in an A/D converter, etc.) affects greater, resulting in a larger number of rise pulses being produced. By this, the compensation coefficient tends to diverge at the time of calculating the error.\nOn the contrary, if the step size parameter μ is set smaller when the transmission signal level is small, a minute error having been detected is canceled, which produces a problem of preventing proper update of the compensation coefficient (refer to FIG. 7A)."}
{"text": "This invention relates to artificial kidney systems; and more particularly, to a proportioningtype dialysis machine for use in such systems.\nIn an artificial kidney system, a dialysis machine delivers dialysis solution under carefully controlled conditions to a dialyzer or artificial kidney. There are two basic types of dialysis machines. One is known as a batch-type machine and the other is a proportioning or continuous feed machine.\nIn both types of machines dialysis concentrate is mixed with water to provide the dialysis solution which is delivered to the dialyzer. Both machines also include means for sterilizing the machine, a heater system for maintaining the dialysis solution at body or physiological temperature (about 37.degree. C.), a device for removing gas and minimizing gas in the dialysis solution, and means for controlling the pressure and flow rate of the dialysis solution in the dialyzer. Furthermore in proportioning machines, provisions are also made for continuously preparing the dialysis solution by mixing or proportioning the water and dialysis concentrate and for controlling the composition of the dialysis solution.\nProportioning machines are generally more compact and smaller in size than the batch-type machine. For this reason, among others, the proportioning machines have experienced increased popularity.\nExisting proportioning machines generally include the required gas removers, heaters, pressure controls, etc., but have been found to be deficient in certain respects. For example, one type of gas removal system, as shown in U.S. Pat. No. 3,738,382, includes a heater for heating the water to a high temperature and a debubbling chamber for removing gas from the heated water at atmospheric pressure. This system does not effectively degass the water and the heating of the water caused dissolved minerals to precipitate and clog passageways within the dialysis machine.\nA second type of gas removal system is shown in U.S. Pat. No. 3,528,550. In this system water is fed to a degassing chamber which is maintained at a pressure below atmospheric pressure by a venturi through which dialysis solution flows. Thus the pressure in the degassing chamber is directly related to the dialysis solution flow rate through the venturi. The venturi only applies a moderate negative pressure to the degassing tank and thus does not effectively degass the water.\nIt is therefore an object of this invention to provide an effective degassing system, and particularly for use in a proportioning dialysis machine.\nIn the degassing system shown in U.S. Pat. No. 3,528,550, the degassing chamber pressure may vary with dialysis solution flow rate, which, in turn, may vary with dialysis conditions, such as patient size, etc. Variations in degassing chamber pressure may affect gas removal.\nIt is therefore an object of this invention to provide a degassing system which functions independently of the dialysis solution flow rates.\nDuring dialysis it is desirable to control the dialysis solution pressure in the dialyzer. However, changes in the dialysis solution flow rate through the dialyzer cause the dialysis solution pressure to vary.\nIt is also an object of this invention to provide means for controlling the dialysis solution pressure in the dialyzer as the dialysis solution flow rate changes.\nThese and other objects of this invention will become apparent from the following description and appended claims."}
{"text": "1. Field of the Invention\nThe present invention relates generally to magnetic devices and magnetic recording media using ferromagnetic films and to devices employing ferroelectric films. The present invention is particularly directed to improvements on the dimensions of the elements and magnetic particles so as to achieve uniform coercive force and coercive field characteristics and to contribute to miniaturization of the devices.\n2. Description of the Related Art\nVarious devices utilizing the giant magnetoresistive (GMR) effect and the tunneling magnetoresistive (TMR) effect have been known in the art. Examples of such devices include magnetic recording elements and magnetic read heads.\nFIG. 1 shows a basic structure of a magnetic random access memory. An example of the magnetic random access memory can be found in Wang et al., IEEE Trans. Magn. 33 (1997), 4498. Referring to FIG. 1, the magnetic random access memory is basically constituted from memory elements which are either GMR elements or TMR elements, word lines, and bit lines, which also function as sense lines. The word lines are orthogonal to the bit lines, and the memory elements are held between the word lines and the bit lines at the intersections thereof. Note that in FIG. 1, reference symbol W denotes the width of each memory element in the direction parallel to the bit lines and reference symbol L denotes the length of each memory element in the direction parallel to the word lines.\nReferring now to FIG. 2, a first end of the memory element is connected to the bit line, and a second end of the memory element is connected via a lead to a logic circuit that selects a memory cell. In the example shown in FIG. 2, a field effect transistor (FET) constituted from a silicon substrate, a drain D, a source S and a gate G is connected to the second end of the memory element via a plug (interconnecting lead). Note that FIG. 2 illustrates an example using a TMR element. The TMR element is constituted from a free layer composed of CoFe, NiFe, or the like, a barrier layer composed of Al2O3 or the like, a reference layer composed of CoFe or the like, a non-magnetic layer composed of Ru or the like, a fixed layer composed of CoFe or the like, and an antiferromagnetic layer composed of PtMn or the like. These layers are arranged in that order when viewed from the bit line side. A GMR element has basically the same multilayer structure as that of the above-described TMR element and only differs from the TMR element in structural details such as the absence of barrier layer, etc.\nThe combination of the anisotropic magnetic field in the soft magnetic free layer and the demagnetization field determined by the size of the free layer defines the magnetic field necessary for rotating the magnetization direction of the free layer, i.e., the coercive force Hc.\nFIG. 3 is an enlarged perspective view of the free layer, the reference layer, and the fixed layer of the memory element. In the drawing, the bold arrow in each layer indicates the magnetization direction of that layer. As shown in the drawing, the x axis extends along the long side of rectangular layers and the y axis extends orthogonal to the x axis. Reference symbol W denotes the width of the memory element in the y direction, and reference symbol L denotes the length of the memory element in the x direction. As shown in FIG. 3, easy axes of the free layer and the fixed layer are substantially parallel to each other. The magnetization direction of the reference layer is antiparallel to those of the free layer and the fixed layer.\nThe magnetization of the fixed layer is fixed by the antiferromagnetic layer. Given, for example, that bit information “1” is represented by the magnetization direction of the free layer being oriented in a direction parallel to the magnetization direction of the reference layer and that bit information “0” is represented by the magnetization direction of the free layer being oriented in a direction antiparallel to the magnetization direction of the reference layer, the magnetization direction of the free layer rotates due to the magnetic field, induced by a bit line current and a word line current, that exceeds the above-described coercive force Hc. Magnetic recording is performed through such a rotation.\nFIGS. 4A to 4C illustrate example structures of shield-type magnetic read heads including read elements each disposed in the gap between a pair of shields (for example, refer to C. Tsang et al., IEEE Trans. Magn. 30 (1994), 3801). The shield-type magnetic read head includes a read element, namely, a GMR element or a TMR element, a lower shield S1, and an upper shield S2. In the drawings, hard magnet layers for controlling magnetic domains, disposed adjacent to the read element, a write head integrally formed above the read head, and so on are omitted for the sake of simplicity of explanation. The GMR or TMR heads shown in FIGS. 4A to 4C are classified into three types according to the direction of the sense current Is. FIG. 4A illustrates, for example, a horizontal current-in-plane (CIP) GMR head in which an electric current flows in the track direction. FIG. 4B illustrates, for example, a vertical CIP GMR head in which an electric current flows in the height direction of the element. FIG. 4C illustrates, for example, a current-perpendicular-to-plane (CPP) GMR or TMR head in which a sense current Is flows in the thickness direction. The view of FIG. 4C is made partially transparent for the sake of simplicity of explanation. In the FIGS. 4A to 4C, arrows in strip-shaped recording media represent the recording magnetization direction.\nFIG. 5 illustrates another example structure including a TMR element. The TMR element is constituted from a free layer, a barrier layer, a reference layer, a nonmagnetic layer, a fixed layer, and an antiferromagnetic layer arranged along the z axis in the drawing. In short, the TMR element in FIG. 5 has substantially the same layer structure as that shown in FIG. 2. However, in FIG. 5, hard magnets #1 and #2 are formed at the two sides of the TMR element, and the nonmagnetic layers NM are formed between the TMR element and the hard magnets #1 and #2. Hard magnet layers for controlling the magnetic domains must be disposed at the two sides of the TMR element to orient the magnetization direction of the free layer in the x axis direction.\nFIG. 6 is an enlarged perspective view of the free layer, the reference layer, and the fixed layer of the above TMR element. The hard magnet layers for controlling the magnetic domains, a nonmagnetic layer, a base layer, and a protection layer are omitted from the drawing. In the drawing, a bold arrow in each layer indicates the magnetization direction of that layer. As shown in FIG. 6, the x axis extends along the long side of the rectangular layers and the y axis extends orthogonal to the x axis. Reference symbol W denotes the width of the element in the y direction, and reference symbol L denotes the length of the element in the x direction. As shown in FIG. 6, whereas the easy axis of the free layer extends substantially in the x-axis direction, the easy axes of the reference layer and the fixed layer are orthogonal to the easy axis of the free layer. The magnetization directions of the reference layer and the fixed layer are antiparallel to each other. The magnetization direction of the fixed layer is pinned by the antiferromagnetic layer.\nFIGS. 7A and 7B are schematic plan views showing the shape of the above-described memory element and the read element. The memory element and read element are formed into a rectangular shape, as shown in FIG. 7A, or into an elliptic shape, as shown in FIG. 7B. In FIG. 7A, reference symbol W denotes the breadth of the element, and reference symbol L denotes the longitudinal length of the element. In FIG. 7B, reference symbol W denotes the length of the short axis and reference symbol L denotes the length of the long axis. In FIG. 7B, the element is formed into an elliptic shape to make the demagnetizing field as uniform as possible inside the element.\nKnown magnetic recording elements and magnetic read elements, however, suffer from the following technical bottlenecks:\n(1) Since the variation in the coercive force Hc among magnetic memory elements is large, practical production of the magnetic memory element is difficult; and\n(2) For magnetic read elements, as the element size is reduced, a decrease in the sensitivity occurs due to the hard magnet layers for controlling magnetic domains, and thus magnetic read heads for higher-density media are difficult to design.\nFirst, regarding point (1) above, in order to put magnetic random access memories having a storage capacity comparable to that of current widespread flash RAMs or DRAMs into practical use, all 106 to 109 magnetic memory elements must exhibit a uniform coercive force Hc. If the variation in Hc is 50% or more, the magnetization rotating current may differ by as much as 200% or more between some elements. In practice, this precludes selective recording.\nAs the size of magnetic memory elements decreases, the ratio of the demagnetizing field in coercive force Hc increases. Since the demagnetizing field is heavily dependent on the size and shape of the element, the variation in the coercive force Hc tends to increase in inverse proportion to the size of the elements. For the purpose of illustration, the dependency of the coercive force Hc on the element size was examined using a magnetic memory element having a free layer made of CoFe 2 nm in thickness. The results are shown in FIGS. 8 and 9.\nIn the graph of FIG. 8, characteristics of square and rectangular elements, i.e., box-shaped elements, as shown in FIG. 7A, when the aspect ratio W:L is varied are shown. The abscissa represents 1/W (unit: 1/μm), and the ordinate represents the coercive force Hc (unit: Oe=103/4π A/m).\nIn the graph of FIG. 9, characteristics of elliptic elements, as shown in FIG. 7B, are shown. The aspect ratio W:L is varied. The abscissa and the ordinate are the same as in FIG. 8.\nIn these graphs of FIGS. 8 and 9, solid lines connecting circular symbols represent the characteristics of an element having an aspect ratio W:L=1:1, long dotted lines connecting square symbols represent the characteristics of an element having an aspect ratio W:L=1:2, and short dotted lines connecting rhombus symbols represent the characteristics of an element having an aspect ratio W:L=1:3. An inclined straight line extending from the origin represents the characteristic according to the theoretical formula: Hc=2πMs(t/W), wherein Ms denotes the intensity of the magnetization, and t denotes the thickness.\nAs shown in these graphs, theoretically, the demagnetizing field is supposed to increase in inverse proportion to the length W of the short side of the element. However, this is not the case in practice. Particularly when the aspect ratio W:L is low, deviation from the theoretical value is significant. Since the number of the magnetic memory elements is large, as described above, it is preferred that the coercive force Hc be constant although some degree of variation in the aspect ratio may be observed.\nSecondly, regarding point (2) above, in order to read higher-density recording media, the size of the read element installed in the magnetic read head must be reduced. Such a reduction in size increases the relative thickness of the domain walls at the ends of the free layer of the read element.\nFIG. 10 is a graph showing the resistance/magnetic field characteristics of read elements each having a NiFe free layer with a track width L of 0.1 μm and 2 nm in thickness. No hard magnet layers for controlling the magnetic domains are provided for the read elements. The abscissa indicates the magnetic field (unit: Oe=103/4π A/m) in the y axis direction in FIG. 6, and the ordinate indicates the resistance R (arbitrary units). In the graph, the solid lines connecting circular symbols represent the characteristics of read elements having an aspect ratio W:L=1:1, the long dotted lines connecting square symbols represent the characteristics of read elements having an aspect ratio W:L=1:2, and the short dotted lines connecting rhombus symbols represent the characteristic of read elements having an aspect ratio of 1:3.\nAs the thickness of the magnetic walls increases, hysteresis appears in the resistance/magnetic field curves of the read elements. In other words, when the aspect ratio is large, the curve around zero magnetic field is a straight line; however, the curve clearly exhibits hysteresis as the aspect ratio decreases.\nThe aspect ratio is preferably as small as is feasibly possible so as to read higher density media. However, read elements must maintain a particular magnetization state when no signals are provided and must respond linearly to an external magnetic field. In other words, read elements with a large coercive force Hc and hysteresis are not preferred. In order to eliminate hysteresis, hard magnet layers for controlling the magnetic domains are provided at the two sides of the element so as to forcibly orient the magnetization direction of the free layer in the track direction by the biasing magnetic field from these hard magnet layers. However, application of the biasing magnetic field from the hard magnet layers inhibits the rotation of the magnetization of the free layer, thereby drastically decreasing the sensitivity to the external magnetic field. In other words, small elements suffer from a decrease in the sensitivity because of the presence of the hard magnet layers.\nAs for the problem of point (1) above, defect-free bulk memories can be manufactured by making the coercive force Hc uniform over all the elements. In contrast, as for the problem of point (2) above, the coercive force Hc needs to be eliminated so as to decrease the biasing magnetic field of the hard magnet layers as much as possible and to manufacture magnetic read heads having high sensitivity. In other words, the object in point (1) is to maintain the coercive force Hc in the hysteresis curve in the easy axis direction of the free layer at a particular level; and the object in point (2) is to reduce the coercive force to zero in the hysteresis curve in the hard axis direction of the free layer. These objects appear contradictory but can be achieved simultaneously if the magnetization of the free layer can rotate simultaneously. The means for achieving these objects is the same.\nThe dependency of the coercive force characteristic on the size and shape of the element can be confirmed through investigation of the magnetization distribution immediately before rotation.\nFIGS. 11A, 11B, 12A, and 12B show example magnetization distributions immediately before the rotation of the magnetization. The distributions are estimated by carrying out a micromagnetics simulation. FIGS. 11A and 11B show the distributions in rectangular elements, and FIGS. 12A and 12B show the distributions in elliptic elements. In the drawings, arrows inside the frames represent magnetization distributions, and streamlines above the frames schematically illustrate the general directions of the magnetization.\nA large rectangular element of, for example, 1/W=1.5 μm−1, has a vortex distribution, as shown in FIG. 11A. A small rectangular element of, for example, 1/W=3.0 μm−1, has a distribution resembling the shape of letter S (hereinafter, the “S distribution”), as shown in FIG. 11B.\nA large elliptic element of, for example, 1/W=1.5 μm−1, has a vortex distribution, as shown in FIG. 12A. A small elliptic element of, for example, 1/W=3.0 μm−1, has a distribution resembling the shape of letter C (hereinafter, the “C distribution”), as shown in FIG. 12B.\nNote that the C and S distributions can be considered as a low-level buckling magnetic wall. Whereas buckling magnetic wall has many undulations, the number of undulations in the C and S distributions is low. The state of the magnetic wall in the element is determined according to the balance between the demagnetizing field energy, the anisotropic energy, the exchange coupling energy, and the Zeeman energy.\nAs the size of the element is reduced, the relative strength of the demagnetizing field energy increases, resulting in changes in magnetization distribution immediately before the rotation of the magnetization. Presumably, in most cases, the vortex distribution has a relatively large coercive force Hc, and the S and C distributions have a relatively small coercive force Hc.\nIt should be noted that the above description can be applied not only to devices using ferromagnetic layers but also to devices using ferroelectric layers. When applied to devices using ferroelectric layers, the demagnetizing field corresponds to the depolarization field, and the coercive force Hc corresponds to the coercive field Ec."}
{"text": "A power conversion device mounted on a vehicle or the like has functions of converting DC power into AC power and supplying the same to a rotating electrical machine and converting AC power from the rotating electrical machine into DC power. The power conversion device has an inverter circuit constituted by a semiconductor element having a switching function. As a circuit body performing power conversion, that is, a power semiconductor module, one with a structure formed by resin-sealing an upper arm circuit and a lower arm circuit, which are constituted by insulating gate bipolar transistors (IGBTs) and diodes, integrally has been known. In a circuit body with this structure, each of the IGBTs and the diodes of the upper and lower arm circuits is mounted on one face of an insulating board. A metal base is arranged on the other face of a pair of insulating boards on which the upper arm circuit or the lower arm circuit is formed.\nThe connection conductors connected to the IGBTs and the IGBTs of the upper and lower arm circuits are mounted so as to form a loop current path on the metal base. In this circuit, when the IGBTs of the upper arm circuit are turned on, the diodes of the lower arm circuit are reverse biased so that the recovery current passes through the upper and lower arm circuits. At this time, an induced current is generated at the metal base. The direction of the magnetic flux generated around this induced current is opposite to the direction of the magnetic flux generated by the recovery current flowing through each conductor plate of the upper and lower arm circuits. Thus, the magnetic fluxes cancel each other, and an inductance of an internal circuit decreases (e.g., see FIG. 9 of PTL 1)."}
{"text": "A nonaqueous secondary battery, which is represented by a lithium ion secondary battery, has a high energy density and is widely used as a main electric power source of a portable electronic equipment, such as a portable phone and a notebook computer. The lithium ion secondary battery is demanded to attain a further high energy density, but has a technical issue on assuring safety.\nA separator plays an important role on assuring safety of a lithium ion secondary battery, and under the current situation, a polyethylene microporous membrane is used since it has a high strength and shutdown function. The shutdown function referred herein means a function of shutting down an electric current by closing the pores of the microporous membrane when the temperature of the battery is increased, and the battery is suppressed from generating heat by the function, thereby preventing the battery from suffering thermal runaway.\nThe energy density of the lithium ion secondary battery is being increased year by year, and for assuring safety, heat resistance is demanded in addition to the shutdown function. However, the shutdown function contradicts the heat resistance since the operation mechanism thereof depends on closure of the pores through melting of polyethylene. There have been proposals on improvement in heat resistance with the molecular weight of polyethylene, the crystalline structure or the like, but sufficient heat resistance has not yet been attained. Such techniques have been proposed that polypropylene is blended or laminated, but under the current situation, these systems fail to attain sufficient heat resistance. Furthermore, for enhancing the heat resistance with the shutdown function attained simultaneously, such techniques have been proposed that heat resistant porous layers are coated on both front and back surfaces of a polyethylene microporous membrane, and nonwoven fabrics containing heat resistant fibers are laminated thereon.\nIt is an important factor of a separator for assuring safety of a nonaqueous secondary battery that the separator has shutdown function and heat resistance, and furthermore, it is also important that the separator has flame retardancy from the standpoint of ignition. The currently available separator for a nonaqueous secondary battery as described above uses a polyethylene microporous membrane in consideration of shutdown characteristics, and there are many techniques for enhancing heat resistance mainly with the polyethylene microporous membrane. Polyethylene is a polymer that is highly combustible, and in consideration of the property, cannot be considered as having high safety.\nSuch a separator has been known that has a polyethylene microporous membrane and a heat resistant porous layer having an oxygen index of 26 or more, which are laminated on each other (see Patent Document 1). However, a polyethylene microporous membrane is still combustible even though it is coated with a layer having a high oxygen index, and it is not effective from the standpoint of flame retardancy.\nSuch a separator has been also known that has a polyethylene microporous membrane and a heat resistant porous layer laminated on each other, in which ceramic powder is mixed in the heat resistant porous layer (see Patent Document 2). In Patent Document 2, the ceramic powder is mixed for the purpose of improving the ion permeability. However, there is no effect in flame retardancy by adding ceramic powder, which is represented by a so-called metallic oxide. Furthermore, the separator has a handling problem, in which an equipment is severely abraded due to the ceramic particles, which are generally hard. In the case where the equipment is abraded, metallic powder and the like are attached to the separator and may cause decrease in capability of the battery.\nIn addition, techniques for imparting flame retardant effect to the separator by adding a flame retarder thereto (see Patent Documents 3 to 6). Patent Document 3 discloses examples of utilizing a halogen flame retarder and barium sulfate in the form of solid particles. Patent Documents 4 to 6 disclose examples of adding a polymer flame retarder to a separator. The proposals contribute to flame retardancy of a separator, but cannot enhance the heat resistance sufficiently, and thus it is difficult to assure safety of a battery.    Patent Document 1: JP-A-2006-269359    Patent Document 2: Japanese Patent No. 3,175,730    Patent Document 3: JP-A-7-272762    Patent Document 4: JP-A-2006-351316    Patent Document 5: JP-A-2005-149881    Patent Document 6: JP-A-2001-210314"}
{"text": "Asymmetric digital subscriber line (ADSL) systems enable data to be transmitted over a pair of metallic twisted pair (usually copper) wires to customer premises. It is thought that the maximum transmission performance that is likely to be obtained with modern variants of ADSL is a download data rate of 24 Mbps and an upload speed of about 3 Mbps. Such data rates are dependent on the length of the metallic twisted pair from the customer premises to the telephone exchange and thus many customers will receive services at significantly lower data rates.\nTo improve data rates optical fibre has been installed into the access network. The greatest data rates are likely to be provided using fibre to the premises (FTTP) networks, such as passive optical networks (PONs), but there is a significant cost involved in providing fibre to customer premises. Fibre to the cabinet (FTTCab) networks are known to provide an attractive solution to providing customers with high data rate services without requiring as much investment as FTTP networks. Typically in FTTCab networks, very high bit-rate digital subscriber line (VDSL) systems are used to provide data rates of 40 Mbps and higher, for both upload and download on the metallic twisted pair cables. It is believed that improvements to VDSL systems may provide data rates in excess of 100 Mbps.\nDSL systems work by utilising the frequencies above those which are used by the conventional telephony signals. In particular, VDSL2 defines three frequency windows for downstream data and 2 frequency windows for upstream data. Each of these windows comprises a number of carriers which have a 4.3125 kHz frequency separation. Each of these carriers will transmit one or more symbols with each of these symbols being used to transmit up to 15 bits of data. During a training process the insertion loss and noise level are determined for each of the carriers such that the signal to noise ratio (SNR) for each carrier can be determined. The training process determines the capacity of the upstream and downstream links in accordance with the SNRs of each of the carriers.\nAccording to a first aspect of the present invention there is provided an apparatus having a first digital subscriber line connection to a first communications network and a second connection to a local area network, the apparatus being configured to, in use, transmit data via the local area network in response to the initiation of a training process for the digital subscriber line connection.\nThe apparatus may comprise a modem or a router. The apparatus may be configured, in use, to transmit data to a device connected to the local area network via a powerline adaptor. Data may be transmitted by the apparatus to a set top box or a router via the powerline adaptor. The apparatus may be further configured to cease transmitting data via the local area network in response to the termination of the training process for the digital subscriber line connection.\nAccording to a second aspect of the present invention there is provided a method of operating a communications network, the method comprising the steps of: a) initiating a training process for a digital subscriber line connection with an apparatus connected to a communications network; and b) in response to the initiation of the training process, transmitting data from the apparatus to a local area network. In step b) the apparatus may transmit data to a further device via a powerline adaptor.\nThe network may comprise the further step of c) ceasing to transmit data to the local area network in response to the termination of the training process for the digital subscriber line connection.\nAccording to a third aspect of the present invention there is provided a tangible data carrier for use in a computing device, the data carrier comprising computer executable code which, in use, performs a method as described above."}
{"text": "The present invention relates generally to telephone message recording and transcribing systems, and in particular to an improved message retrieval system for use with a telephone answering system which includes a plurality of telephone answering machines and a plurality of message memory units coupled to the answering machines.\nTelephone answering systems are routinely used to enable operating personnel to service a large number of callers. Such systems include a plurality of answering machines coupled to telephone lines. These answering machines deliver prerecorded messages after seizing a line and record caller responses on a message memory unit such as an endless loop tape recorder. Operators then transcribe the recorded messages on the message memory units and take appropriate action. For example, an operator may call back a caller who has left a telephone number and a request.\nU.S. Pat. No. 4,338,494 discloses a microprocessor-controlled telephone call inventorying and sequencing system. This system records priority information (such as the time of recording of a message) along with the stored messages in the message memory unit. The microprocessor uses this stored priority information to select which of the message memory units is to be interconnected with a transcribing station requesting transcription. The selection criteria are chosen to prevent stored messages on any particular message memory unit from being delayed excessively in transcription. The system disclosed in the above identified patent operates automatically to disconnect a message memory unit from a transcribing station if the message memory unit has not been transcribed from in a predetermined period, such as 60 seconds. In this way, the system prevents an inactive transcribing station from monopolizing a message memory unit.\nOur previous U.S. Pat. No. 4,150,255 discloses a conversational telephone call distributor which utilizes a manually controlled distribution panel to interconnect operator stations with selected message recorders.\nIn spite of the many advantages offered by the systems described above, a need presently exists for an improved message retrieval system which provides improved protection against operator abuse of the system or the callers, improved flexibility by which an operator can both transcribe recorded messages and respond to live telephone calls, and which takes positive action to prevent stored messages from being delayed excessively in transcription."}
{"text": "A DC/DC voltage regulator is operative to maintain a level output voltage despite variations in power supply voltage or current drawn by a load. As one example, a step down converter may take a relatively unstable input voltage and maintain a desired level output voltage that is nominally lower than the input voltage. Specifically, a step down converter may receive as an input a voltage in the range of 2-10 volts and output a level voltage of (for example) 1.8 volts.\nMany portable devices require a steady voltage supply such as that provided by a DC/DC voltage regulator. Further, the advent of portable electronics and the need for longer battery life requires new types of voltage regulators. These DC/DC voltage regulators need to be efficient while operating in both low current and high current load conditions. For example, handheld electronics such as PDA's and cell phones now require high efficiency at varying loads (such as standby and active modes) to extend battery life. The standby mode requires a very low amount of current to operate. Only critical systems and volatile memory need to be powered to constantly refresh and maintain the data in the device. Because of these requirements, new voltage regulator schemes have been developed that are very efficient at all current levels.\nFor example, a pulse width modulated (PWM) switch-mode regulator is an efficient regulation scheme during heavy loads. It offers high efficiency, low output voltage ripple, good line and load regulation. However at light loads the PWM regulator has poor efficiency.\nAt light loads, a pulsed frequency modulation (PFM) switch-mode regulator is commonly used due to its high efficiency. However the large output voltage ripple, poor line and load regulation inherent to PFM precludes its use in many systems. Thus, a low quiescent current LDO regulator is desirable in these systems. The LDO regulator offers relatively good light load efficiency, low output voltage ripple, and good line and load regulation, but at heavy loads the efficiency is far below that of switching regulators.\nCombining a PWM switch-mode regulator and a linear LDO regulator in parallel offers the high efficiency and good output voltage regulation required by many portable battery powered systems. Switch-mode and linear regulators have been paralleled in the prior art, but the regulation voltages have been slightly different, and the control scheme of the regulators are very basic. These types of systems do not have the optimal efficiency and the output voltage regulates at two different voltages making the load regulation poor."}
{"text": "The Personal Computer Memory Card International Association (PCMCIA) is a memory card standard for the interface between a computer and a computer peripheral product. A PCMCIA card (briefly referred to PC card) is an integrated circuit (IC) card that conforms to the PCMCIA standard, and a PCMCIA slot is a slot that conforms this standard. PC cards are widely used for connecting computers and computer peripheral devices. Especially for notebook computers, external peripheral devices are required due to the small size of a notebook computer, and PC cards therefore become an optimum interface. Presently, PC cards mainly act as data transmitting or local network cards for notebook computers, in order to eliminate the user's trouble for carrying heavy peripheral devices.\nIn order to provide more a superior transporting interface, the PCMCIA developed an express card which is lighter, faster, thinner, easier to use, and applies for more extensive I/O mode than traditional PC cards. The express card has two standards, one is express card/34 having a 34 mm width (small card) and the other is express card/54 having a 54 mm width (big card).\nIn general, when inserting into a card connector, the express card is positioned in the slot of the card connector and pushed inward to the top of the slot. During insertion, the user has to aim the express card at the opening of card connector carefully thus making a quick insertion difficult. FIG. 7, FIG. 8, and FIG. 9 respectively show a top view of a known card connector, a top view of a known top board, and the cross-section view of A-A line in FIG. 8. When inserting the card (the direction shown by the arrow “a” in FIG. 7 and FIG. 8), if the user can not precisely aim the card at the right position of cover 2 of the card connector (the center of the cover 2), the cover 2 will be undesirably urged and bent in the card-inserting direction “a”, and can be easily damaged. For example the cover 2 can be damaged such that it separates from the card connector."}
{"text": "This invention relates to the aerially transporting of loads between sites and more particularly, to the transporting of such loads using a lift balloon.\nThe concept of aerially transporting loads such as logs with lift balloons is well known. Balloons provide a great amount of lift at relatively low cost and can be maneuvered to hover over a loading or discharging site.\nA recent example of a method for aerially transporting loads between loading and discharging sites by balloons is disclosed in U.S. Pat. No. 4,055,316 to Chipper et al. There, the lift balloon is guided between the sites by an aerial cable. The balloon contains a liquid ballast container that can be filled with ballast to maintain the proper buoyancy of the balloon. At the loading site, liquid ballast is ejected in an amount corresponding to the weight of the load to be suspended from the balloon. Upon return to the load discharging site, liquid ballast is replaced in the ballast container and the load is removed from the balloon.\nOne drawback of Chipper et al. and other conventional balloons is the aerodynamic drag their pear shape creates. Specially shaped balloons may be formed to reduce this drag, such as shown in U.S. Pat. No. 3,369,673 to Mosher, but these balloons cost considerably more than conventional balloons.\nAnother drawback is the difficulty of maintaining the proper buoyancy of the balloon with the added weight of the load. Chipper et al. compensates by ejecting a corresponding amount of ballast. But this ballast is then lost and must be replaced by additional ballast. However, replacement may be impractical under conditions where water is not readily available, such as in remote logging areas."}
{"text": "1. Field of the Invention\nThis invention relates to a dry etching method adapted for fine processing of manufacturing a semiconductor device, and particularly to a method for preventing regression of a resist mask formed on an SiON based antireflection film so as to improve anisotropy.\n2. Description of Related Art\nIn order to realize large scale integration of semiconductor devices, the minimum processing size of the circuit pattern formation has been rapidly diminished. For instance, the minimum processing size of the 16M DRAM of approximately 0.5 .mu.m (half micron), the minimum processing size of the 64M DRAM of 0.35 .mu.m (sub-half micron), and the minimum processing size of the 256M DRAM of 0.25 .mu.m (quarter micron) are required.\nThis increasingly fine processing depends largely upon a technique of photolithography to form a mask pattern. Visible to near ultraviolet rays, such as g rays having a wavelength of 436 nm or i rays having a wavelength of 365 nm, of a high pressure mercury lamp are used for the current 0.5-.mu.m class processing, and far ultraviolet rays, such as KrF excimer laser lights having a wavelength of 248 nm, are used for 0.35 to 0.25-.mu.m class processing. In the photolithography technique for forming a fine mask with a ray width of not greater than 0.4 .mu.m, an antireflection film to weaken a reflected light from an underlying material layer is substantially required for preventing reduction in contrast and resolution due to halation and standing wave effect.\nAs the component material of the antireflection film, amorphous silicon, TiN and TiON are conventionally used. However, since it has been shown that SiON (silicon oxide nitride) exhibits satisfactory optical properties in the far ultraviolet region, application of SiON to the excimer laser lithography is proposed. It is exemplified by a process of fine gate processing with an SiON film restraining the reflectivity of a W (tungsten)--polycide film or an Al (aluminum) based material film.\nMeanwhile, after the patterning of the resist mask by such photolithography is finished, the antireflection film is etched in the subsequent etching process.\nIn this case, such a problem is now being apparent that the anisotropic shape of the underlying material layer may be deteriorated by oxygen discharged from SiON in the etching process, particularly in overetching. This problem is explained with reference to FIGS. 1 to 4.\nFIG. 1 shows a state of a wafer prior to the etching, in which a gate SiO.sub.x film 22, a W-polycide film 25 and an SiON antireflection film 26 are sequentially stacked on an Si substrate 21, with a resist mask 27 patterned in a predetermined shape being formed thereon. The W-polycide film 25 is composed of, from the bottom, a polysilicon layer 23 containing impurities and a WSi.sub.x (tungsten silicide) layer 24 which are sequentially stacked.\nIf the W-polycide film 25 is etched using a Cl.sub.2 /O.sub.2 mixed gas, the etching is promoted by a formation of etching reaction products, such as SiCl.sub.x and WClO.sub.x. On the other hand, a carbon based polymer derived from decomposition products provided by forward sputtering of the resist mask is deposited to form a sidewall protection film 28 on the sidewall surface of the pattern. If the wafer temperature is sufficiently low, SiCl.sub.x of relatively low vapor pressure among the etching decomposition products can be a component of the sidewall protection film 28.\nAs a result, a gate electrode 25a of anisotropic shape is formed at the end of just etching, as shown in FIG. 2. In FIG. 2, materials after the etching are denoted by their respective original numerals plus subscripts \"a\".\nHowever, if the overetching follows the just etching, regression of the edge of the resist mask 27 causes the SiON antireflection film 26a to have its end surface tapered to be easily exposed, as shown in FIG. 3. SiON, having an element composition ratio of approximately Si:O:N=2:1:1, is richer in Si than SiO.sub.2 is. Consequently, SiON has low durability to a Cl based plasma, and easily discharges active O* when its exposed end surface is etched. Then, O* removes the sidewall protection film 28 in the form of CO.sub.x, to lower sidewall protection effects. In addition, since the W-polycide film 25 to be etched is reduced in the overetching, a relatively excessive amount of O* is present in the etching gas.\nAs a result, a gate electrode 25b having an undercut is formed, as shown in FIG. 4. The material layers having the undercut denoted by their respective original numerals plus subscripts \"b\". The undercut is generated most conspicuously in the WSi.sub.x layer 24b. Since O* sputtered out from the end surface of the SiON antireflection film removes W atoms in the form of WClO.sub.x, the etchrate in the WSi.sub.x 24a is increased.\nAs the anisotropic shape of the gate electrode is thus deteriorated, serious problems rise, such as, the metallization resistance falling off the designed value and difficulty in forming the sidewall to attain an LDD structure.\nThe deterioration in the anisotropic shape in the overetching is not limited to the above-described SiON antireflection film, but is a phenomenon which may be generated in cases where an antireflection film capable of easily discharging oxygen is used and where conductive material layers of Al based metallization and the like other than the W-polycide film are used as etching targets."}
{"text": "Transistors are used for many purposes. When used as amplifiers, high gain is often desirable. However, it can be desirable for such a transistor to be substantially linear as non-linearity can introduce distortion, and it is typically desirable that the gain of a transistor should not change in response to operating conditions such as collector-emitter voltage Vce. This is generally indicated by a transistor exhibiting a relatively large “Early” voltage.\nThe Early effect (named after its discoverer James Early) describes how an effective width of a base region of a bipolar transistor changes with collector-emitter voltage. The width of the base affects the gain of the transistor. Consequently the gain of a transistor can vary with the instantaneous amplitude of a signal. If the signal was a sinusoid then the gain applied to the peak of the sinusoid would be different to the gain applied to the trough of the sinusoid, which gives rise to harmonic distortion"}
{"text": "In the prior art, all air conditioning units for aircraft, such as helicopters and small airplanes, are air conditioned with permanently installed air conditioners. Generally these units are mounted to the mainframe and may be divided into several different modules. The various components or modules are connected via permanent hoses and electrical connections. Because of the various components and connections, these prior art systems are heavier than a single modular portable unit. Also, because these units are installed permanently, the aircraft must carry the additional weight of the air conditioner whether it is being used or not.\nHowever, a portable unit is difficult to use in small aircraft because ducts cannot be easily installed to remove waist heat from the condenser. Generally, even if portable or semi-portable air conditioners are used the ducting must be permanent and in some instances, such as the BAK-109 used on the Beech 18 and DC-3, the airframe was permanently altered to allow for the ducting. Thus, these prior art air conditioners had to be installed when the aircraft was manufactured or as a retrofit in the factory.\nIt would be highly advantageous, therefore, to remedy the foregoing and other deficiencies inherent in the prior art.\nAccordingly, it is an object of the present invention to provide a new and improved portable vapor cycle air conditioning unit for small aircraft.\nAnother object of the invention is to provide a new and improved vapor cycle air conditioning unit that is easy to install and remove.\nAnother object of the invention is to provide a new and improved vapor cycle air conditioning unit that does not require any permanent alterations of the aircraft.\nAnother object of the invention is to provide a new and improved vapor cycle air conditioning unit in which normal condensate within the air conditioning unit is eliminated without requiring any permanent alterations of the aircraft.\nStill another object of the invention is to provide a new and improved vapor cycle air conditioning unit that requires a very small space."}
{"text": "Many attempts have been made to utilize simultaneously the allied foot and arm powers while the arms also steer the bicycle. In the prior art the steering mechanisms pivot with the front wheel fork while at the same time imparting power to assist turning the wheels. Because of this dual function, inefficencies result. For a person to crank efficiently the hand crank rotational axis should be substantially parallel to the person's shoulder line, while the imaginary longitudinal center plane of his body lies in coincidence to a centrally located vertically oriented plane of the hand crank. In past devices the structures are such that, when their bars or handles are turned to the right or left, the above-referred to planes will automatically produce an obtuse included angle. In these instances the person cannot transmit an effective manual crank power because the anatomy of the human arms simply cannot function effectively in this position."}
{"text": "Multi-layered heteostructures are employed to implement devices for a number of applications. These applications include, but are not limited to, optoelectronic components (e.g. PIN junction or multi-quantum well). The functionality of these multi-layered heteostructures are typically built from layer to layer in a vertical direction, using different semiconductor materials. Further, the multi-layered heteostructures are vertically etched leading to the exposure of their sidewalls, and polymer is spun to seal the sidewalls. To facilitate provision of a contact to one of these devices, the polymer may be etched back to expose the top semiconductor layer, to allow a metal contact to be deposited thereon. Alternatively, a vertical via may be etched to open the polymer to facilitate contact between the top semiconductor layer and the metal contact.\nHowever, both practices have disadvantages. In particular, the former practice may not be able to clear the top semiconductor layer without exposing the sidewalls of some of the device layers underneath the top layer. Whereas, the latter practice is difficult and complicated, especially in the smaller than micro scale, e.g. at nanometer scale. As at the nanometer scale, not only alignment of the via mask becomes very difficult, making of the via mask in and of itself becomes almost impossible, due to current sub-micro lithography is unable to accurately resolve nanometer via printing. Also, at nanometer scale, the via approach will not allow the full use of the available area of the top semiconductor layer because a typical via approach requires some margin so the via must be smaller than the device. Even if the first practice is able to open the whole area of the top device layer, at micrometer or nanometer scale, the top semiconductor area may not be sufficiently large to provide a desired low contact resist interconnect (as resist is inversely proportional to the contact area)."}
{"text": "This invention relates to portable, power-driven, cutting tool guides and gauges, and more particularly to a power saw guide and gauge for crosscut operations on wooden boards and planks.\nA very large percentage of cuts with a hand power saw in construction work involve cutting boards and planks to proper length with a square, crosscut. The cuts are generally repetitious in that many boards or planks must be cut to the same length.\nA problem typically encountered with many available portable cutting tool guides is their need for mechanical clamping devices. Such devices often are not suitable because they are cumbersome and require excessive amounts of time to set up, thereby reducing the saw user's efficiency. Those guides that do not require clamping devices usually do not have a gauge for determining board length, and a separate length measuring operation must take place before a crosscut is made, thereby reducing the saw user's efficiency.\nThe few guides which combine a gauge means, usually lack the ruggedness needed for outdoor construction work and are more readily applicable to indoor or occasional use. None lend themselves to quick and repetitious use."}
{"text": "Manufacturers of personal care and cosmetic products typically use a variety of emulsions. This requires a certain supply chain and manufacturing complexity. Using a common formulation to create a variety of products ranging from lotions to creams would greatly simplify the manufacturing supply chain and reduce overall manufacturing costs. However, there are limitations to the extent which a formulation can be used in a variety of products because of instability.\nWhat is needed is a common base emulsion that can be used to make a variety of cosmetics. It is desirable to have a common base emulsion that is sprayable. Further, it is desirable to have a common base emulsion that is largely anhydrous, yet stable."}
{"text": "1. Field of the Invention\nThis invention relates to a collapsible work horse having first and second pairs of legs pivotally mounted to a support beam to move from an extended or working position to a storage and/or transporting, e.g., collapsed, position, and a locking arrangement to lock the legs in the extended position and, more particularly, to a collapsible work horse having the legs secured in the extended position by a plunger mounted in each of the legs and biased into a hole in the support beam. The invention further relates to a work station having one or more work horses for supporting a shaping tool and for supporting the pieces to be shaped.\n2. Discussion of the Technical Problems\nIn general, work horses, also known as sawhorses or trestles, include a first pair of legs secured to one side of a support beam and a second pair of legs secured to an opposite side of the support beam. The legs can be fixedly secured to the support beam using fasteners, e.g. but not limited to, nails, screws, and/or nut and bolt arrangements, or detachably secured to the support beam using clamps. In general, the clamps include a pair of elongated members pivotally mounted together such that moving one end of the members away from one another moves the opposite ends of the members toward one another against the support beam. In another arrangement, the legs are secured by pivotally attaching the legs to the support beam as taught in U.S. Pat. No. 3,951,233 (hereinafter also referred to as “U.S. Pat. No. '233”).\nAlthough the presently available work horse designs are acceptable for their intended use, they have drawbacks. More particularly, work horses that have the legs and support beam fixedly secured together are usually moved and/or stored in the assembled state, which results in wasted unused space. The work horses that have the legs detachably secured to the support beam reduces the amount of unused space required for storage but requires disassembling the work horse, keeping track of the disassembled parts, and assembling the parts to use the work horse.\nThe collapsible work horse of U.S. Pat. No. '233 eliminates many of the problems discussed above; however, the work horse of U.S. Pat. No. '233 has limitations. More particularly, the extended legs of the work horse disclosed in U.S. Pat. No. '233 are maintained in the extended position by a constant frictional force applied to the pivot point of the legs. The frictional force is applied by tightening the bolt at the pivot point. For a detailed discussion of the arrangement to maintain the legs in the extended position, reference can be made to Patent '233.\nAs can be appreciated, tightening bolts to secure the legs in the extended position requires the use of the tool to tighten the bolts to secure the legs in the extended position and to loosen the bolts to move the legs to the collapsed position. It can be appreciated by those skilled in the art that it would be advantageous to provide a work horse that has legs that can be moved between the extended position and the collapsed position and does not have the drawbacks and/or limitations of the presently available work horses."}
{"text": "Wellbores may be drilled into a surface location or seabed for a variety of exploratory or extraction purposes. For example, a wellbore may be drilled to access fluids, such as liquid and gaseous hydrocarbons, stored in subterranean formations and to extract the fluids from the formations. Wellbores used to produce or extract fluids may be lined with casing around the walls of the wellbore. A variety of drilling methods may be utilized depending partly on the characteristics of the formation through which the wellbore is drilled.\nThe drilling system may drill a wellbore or other borehole through a variety of formations. The formation may include geologic formations ranging from unconsolidated material to rock formations such as granite, basalt, or metamorphic formations. The drilling system may include a drill bit with a plurality of cutting elements located on the bit to loosen material from the formation to create the wellbore. The cutting elements may include a cutting edge or surface on that is sufficiently durable to penetrate through the formation and maintain desirable uptime of the drilling system.\nHarder formations (i.e., geologic formations including harder rocks or other materials) increase wear on a drill bit and the cutting elements mounted on the drill bit compared to softer formations. The increased wear in harder formations increases the risk of failure of a cutting element or the drill bit and, therefore, increases the risk of damage to the drilling system. The increased wear in harder formations reduces the operational lifetime of a cutting element and drill bit, which in-turn increases the time and cost involved in retrieving the drill bit from the wellbore, replacing or repairing the drill bit, and tripping the drill bit back into the wellbore."}
{"text": "The present invention deals with a method of and a device for producing heat/cold with electric energy. The present invention can be used in refrigerators and air-conditioners of high efficiency as well as for heating spaces by taking energy from outer areas having deep cold.\nCold producing machines usually produce cold from a supplied energy, such as mechanical, thermal or electrical energy. In known cooling devices, a working gas is Freon which inevitably escapes to the atmosphere. When refrigerators or air-conditioners are utilized in great numbers, the Freon leaks lead to substantial damages to the ecology. It is therefore extremely undesirable to use such refrigerators and air-conditioners.\nOn the other hand, semiconductor elements are known, in which when electric current is supplied through the semiconductor element, one end of the element becomes hot while the other end of the element becomes cold. Such an element usually includes a negative semiconductor branch and a positive semiconductor branch. The substantial disadvantage of such elements is their high thermal conductivity since they are formed as monolithic semiconductor elements. A direct flow of heat from the hot end to the cold end occurs in these elements, which reduces efficiency of the system and makes it unsuitable for practical use."}
{"text": "Wear and laundering of fabric articles, and particularly white fabric articles, can result in a discoloration from the original fabric color. For example, white fabrics which are repeatedly laundered can exhibit a yellowing in color appearance which causes the fabric to look older and worn. To overcome the undesirable yellowing of white fabrics, and similar discoloration of other light colored fabrics, some laundry detergent products include a hueing or bluing dye which attaches to fabric during the laundry wash and/or rinse cycle.\nHowever, after repeated laundering of fabric with detergent containing bluing dye, the bluing dye tends to accumulate on the fabric, giving the fabric a bluish tint. Such repeated laundering of white fabric articles tends to give the articles a blue, rather than white, appearance. To combat this accumulation of bluing dyes on fabric, chlorine treatments have been developed. While the chlorine treatment is effective to remove accumulated bluing dyes, the chlorine treatment is an additional and often inconvenient step in the laundry process. Additionally, chlorine treatment involves increased laundering costs and is harsh on fabrics and therefore undesirably contributes to increased fabric degradation. Accordingly, a need exists for improved laundry detergents which can counter the undesirable yellowing of white fabrics, and similar discoloration of other light colored fabrics."}
{"text": "Pigments which are inferior in heat resistance, weatherability, light fastness and chemical stability (resistance to acids, alkalis and other chemicals) can be improved by coating the surface thereof with a physically and chemically stable metal oxide or treating the surface thereof with a coupling agent.\nFor instance, cadmium pigments undergo discoloration when subjected to a temperature 700.degree. C. or higher and their brilliant color is lost, because cadmium sulfide, their principle, is converted to cadmium sulfate or cadmium oxide. Therefore, cadmium pigments cannot be used for products such as ceramic tiles which are subjected to high temperatures.\nOf the chromium yellow pigments, the principle of which is lead chromate, and chromium orange are inferior in heat resistance and acid resistance and lemon yellow is inferior in heat resistance and resistance to alkali and they all suffer discoloration. Molybdenum red, which essentially consists of lead chromate and lead molybdate, is inferior in heat resistance, acid resistance, alkali resistance, weatherability and light fastness, and suffers discoloration.\nCadmium yellow (pale light), ultramarine, chromium green, lithopone, Bordeaux-10B, etc. are inferior in acid resistance, and prussian blue, chromium green, brilliant carmine 6B, lake red C, Bordeaux-10B, rhodamine lake Y, etc. are inferior in alkali resistance.\nAs methods for surface treatment of pigment for the purpose of improving heat resistance, weatherability, light fastness, acid resistance, alkali resistance, etc., the following are known; (a) To form a coating of silica on the surface of pigments by addition of a silicate salt when the pigments are formed by precipitation. (b) To form a coating by adding a metallic salt which forms a water-insoluble metal hydroxide or oxide upon neutralization or a metal salt pair which forms a water-insoluble salt by double decomposition reaction in a dispersion of a pigment. (c) To modify the surface of pigments using a titanium coupler or silane coupler.\nIn addition to the above, there is known a method for preventing oxidation decomposition of pigments for ceramic tiles comprising mixing a stable material such as zirconium silicate which is not attacked by glaze into a pigment and firing the mixture.\nIn the methods (a) and (b), precipitate of silica or a water-insoluble metal salt is formed on the surface of pigments. Formation of precipitate is usually influenced by temperature, pH, etc. and adjustment of many factors are required, and yet it is very difficult to form dense and homogeneous coating. The method (c) is effective for improving weatherability, light fastness and dispersibility of pigments. However, the couplers are mere surface-modifiers for pigments and do not form a dense coating on the surface of pigments, and therefore, they are of little use for improving heat resistance, acid resistance, alkali resistance, etc.\nWhen a coating for improving heat resistance, light fastness, acid resistance and alkali resistance of pigment is formed, formation of the coating is usually carried out under the heated condition, since the reactivity is low at room temperature or lower temperatures. By the reaction at an elevated temperature, no dense coating film is formed. The reason is surmised to be that metal hydroxide molecules formed by rapid hydrolysis attach to the already formed suspended metal hydroxide or oxide particles in preference to the surface of pigments, meaning that the percentage of the molecules depositing on the pigment surface is low.\nWe have found that if a hydrophilic pigment is contacted with a metal alkoxide such as alkyl silicate in an at least partially water-miscible organic solvent, a dense and uniform coating is formed at relatively lower temperatures. Further, we found that if a secondary coating is formed on the thus formed primary coating by repeating the above procedure or by the known precipitation process, a firm coating is formed and heat-resistance, weatherability, light fastness, chemical stability such as acid resistance and alkali resistance, etc. of the pigment is further improved."}
{"text": "1. Technical Field\nThe present invention relates to processes for determining torque output and controlling power impact tools. The invention also relates to a mechanical impact wrench having electronic control.\n2. Related Art\nIn the related art, control of power impact tools has been accomplished by directly monitoring the torque of impacts of the tool. For instance, in U.S. Pat. Nos. 5,366,026 and 5,715,894 to Maruyama et al., incorporated herein by reference, controlled impact tightening apparatuses are disclosed in which complex processes involving direct torque measurement are used. Direct torque measurement involves the measurement of the force component of torsional stress, as exhibited by a magnetic field about a tool output shaft, at the point in time of impact. From this force component, related art devices directly determine the torque applied during the impact, i.e., torque T=force F times length of torque arm r. As exemplified by FIG. 10 of U.S. Pat. No. 5,366,026, however, torque measurements fluctuate, even after a large number of impacts are applied. This phenomena is caused by the inconsistent nature of the force component of the impact. In particular, some devices measure torque at a given point in time, such that the torque measured is based on whatever force is being applied at that point in time. In other cases, the force is monitored as it rises, and is measured for peak at a point in time at which a force decrease is detected. In either case outlined above, the force may not be the peak force and, hence, the peak torque derived may not be accurate.\nTo rectify this problem, related art devices use weighting factors, or peak and/or low pass filtering of torque peak measurement, and/or assume, even though it is not the case, a constant driving force from the motor. For instance, in U.S. Pat. No. 5,366,026, torque measurements are used to calculate a clamping force based on the peak value of a pulsatory torque and an increasing coefficient that represents an increasing rate of a clamping force applied. Unfortunately, torque measurement accuracy remains diminished. Accordingly, there exists a need for better processes of operating power impact tools and, in particular mechanical impact tools (i.e., those with mechanical impact transmission mechanisms), with greater accuracy of torque measurement. There also exists a need for more accurate torque measurement.\nAnother shortcoming of the related art is the lack of an electronic control in a mechanical impact wrench."}
{"text": "Most modern vehicles include inflatable restraint apparatus having deployable airbags positioned in many locations throughout an automotive vehicle. Generally, an interior panel includes a deployment door formed into the panel which is designed to break free upon deployment of the airbag.\nA primary aim of the airbag assembly is to control the opening of the deployment door to avoid break explosion and the possibility of flying parts. Clean deployment is achieved in some airbag assemblies by providing a deployment door with a seam, meaning the door is not physically interconnected with the surrounding interior panel. Unfortunately, the seam in the interior panel around the deployment door may not be visually appealing. Thus, in other airbag assemblies clean deployment is provided, in part, by a “seamless” deployment door having aggressive pre-weakening of the outline of the door (typically by laser scoring, mechanical scoring, etc.) by cutting the material or creating perforations.\nWhile this weakening of seamless deployment doors is typically done on the underside of the panel, there still exists a potential for creating blemishes or other disturbances on the exposed class “A” surface. Furthermore, there is a potential to have an uneven break or tear in the deployment door since the plastic is not completely cut through. Accordingly, there exists a need to provide a seamless deployment door which quickly and reliably breaks free from the surrounding panel while eliminating the potential for surface blemishes due to pre-weakening of the door outline."}
{"text": "The structure of the intervertebral disc disposed between the cervical bones in the human spine comprises a peripheral fibrous shroud (the annulus) which circumscribes a spheroid of flexibly deformable material (the nucleus). The nucleus comprises a hydrophilic, elastomeric cartilaginous substance that cushions and supports the separation between the bones while also permitting articulation of the two vertebral bones relative to one another to the extent such articulation is allowed by the other soft tissue and bony structures surrounding the disc. The additional bony structures that define pathways of motion in various modes include the posterior joints (the facets) and the lateral intervertebral joints (the unco-vertebral joints). Soft tissue components, such as ligaments and tendons, constrain the overall segmental motion as well.\nTraumatic, genetic, and long term wearing phenomena contribute to the degeneration of the nucleus in the human spine. This degeneration of this critical disc material, from the hydrated, elastomeric material that supports the separation and flexibility of the vertebral bones, to a flattened and inflexible state, has profound effects on the mobility (instability and limited ranges of appropriate motion) of the segment, and can cause significant pain to the individual suffering from the condition. Although the specific causes of pain in patients suffering from degenerative disc disease of the cervical spine have not been definitively established, it has been recognized that pain may be the result of neurological implications (nerve fibers being compressed) and/or the subsequent degeneration of the surrounding tissues (the arthritic degeneration of the facet joints) as a result of their being overloaded.\nTraditionally, the treatment of choice for physicians caring for patients who suffer from significant degeneration of the cervical intervertebral disc is to remove some, or all, of the damaged disc. In instances in which a sufficient portion of the intervertebral disc material is removed, or in which much of the necessary spacing between the vertebrae has been lost (significant subsidence), restoration of the intervertebral separation is required.\nUnfortunately, until the advent of spine arthroplasty devices, the only methods known to surgeons to maintain the necessary disc height necessitated the immobilization of the segment. Immobilization is generally achieved by attaching metal plates to the anterior or posterior elements of the cervical spine, and the insertion of some osteoconductive material (autograft, allograft, or other porous material) between the adjacent vertebrae of the segment. This immobilization and insertion of osteoconductive material has been utilized in pursuit of a fusion of the bones, which is a procedure carried out on tens of thousands of pain suffering patients per year.\nThis sacrifice of mobility at the immobilized, or fused, segment, however, is not without consequences. It was traditionally held that the patient's surrounding joint segments would accommodate any additional articulation demanded of them during normal motion by virtue of the fused segment's immobility. While this is true over the short-term (provided only one, or at most two, segments have been fused), the effects of this increased range of articulation demanded of these adjacent segments has recently become a concern. Specifically, an increase in the frequency of returning patients who suffer from degeneration at adjacent levels has been reported.\nWhether this increase in adjacent level deterioration is truly associated with rigid fusion, or if it is simply a matter of the individual patient's predisposition to degeneration is unknown. Either way, however, it is clear that a progressive fusion of a long sequence of vertebrae is undesirable from the perspective of the patient's quality of life as well as from the perspective of pushing a patient to undergo multiple operative procedures.\nWhile spine arthroplasty has been developing in theory over the past several decades, and has even seen a number of early attempts in the lumbar spine show promising results, it is only recently that arthoplasty of the spine has become a truly realizable promise. The field of spine arthroplasty has several classes of devices. The most popular among these are: (a) the nucleus replacements, which are characterized by a flexible container filled with an elastomeric material that can mimic the healthy nucleus; and (b) the total disc replacements, which are designed with rigid endplates which house a mechanical articulating structure that attempts to mimic and promote the healthy segmental motion.\nAmong these solutions, the total disc replacements have begun to be regarded as the most probable long-term treatments for patients having moderate to severe lumbar disc degeneration. In the cervical spine, it is likely that these mechanical solutions will also become the treatment of choice. At present, there are two devices being tested clinically in humans for the indication of cervical disc degeneration. The first of these is the Bryan disc, disclosed in part in U.S. Pat. No. 6,001,130. The Bryan disc is comprised of a resilient nucleus body disposed in between concaval-covex upper and lower elements that retain the nucleus between adjacent vertebral bodies in the spine. The concaval-convex elements are L-shaped supports that have anterior wings that accept bones screws for securing to the adjacent vertebral bodies.\nThe second of these devices being clinically tested is the Bristol disc, disclosed substantially in U.S. Pat. No. 6,113,637. The Bristol disc is comprised of two L-shaped elements, with corresponding ones of the legs of each element being interposed between the vertebrae and in opposition to one another. The other of the two legs are disposed outside of the intervertebral space and include screw holes through which the elements may be secured to the corresponding vertebra; the superior element being secured to the upper vertebral body and the inferior element being attached to the lower vertebral body. The opposing portions of each of the elements comprise the articulating surfaces that include an elliptical channel formed in the lower element and a convex hemispherical structure disposed in the channel.\nAs is evident from the above descriptions, the centers of rotation for both of these devices, which are being clinically tested in human subjects, is disposed at some point in the disc space. More particularly with respect to the Bryan disc, the center of rotation is maintained at a central portion of the nucleus, and hence in the center of the disc space. The Bristol disc, as a function of its elongated channel (its elongated axis being oriented along the anterior to posterior direction), has a moving center of rotation which is, at all times maintained within the disc space at the rotational center of the hemispherical ball (near the top of the upper element).\nIt is a principal object of the present invention to provide a cervical disc replacement device that supports, preserves, and maintains proper disc height.\nIt is an equally important object of the present invention to provide a cervical disc replacement that supports, preserves, and promotes proper anatomical motion.\nOther objects of the invention not explicitly stated will be set forth and will be more clearly understood in conjunction with the descriptions of the embodiments disclosed hereafter."}
{"text": "1. Field of the Invention\nThe present invention relates to a waterproof part for a feedhorn, and more particularly, to a waterproof part having double waterproof units for a feedhorn to ensure return loss and useful bandwidth.\n2. Description of the Prior Art\nA feedhorn, which is also known as low-noise block converter, for a satellite antenna is disposed on a focus of a dish reflector of the satellite antenna. The feedhorn is used for receiving radio signals reflected via the dish reflector from a satellite or transmitting radio signals to the satellite. The satellite antenna is usually installed at an outdoor location such as a roof or an exterior wall of a building to ensure communication quality against signal blocking.\nThe feedhorn is normally equipped with a waterproof part made of insulation materials to prevent rain water from dripping into the feedhorn. During signal transmission, the satellite signals encounter insertion loss and part of the satellite signals are attenuated when the satellite signals pass through the waterproof part. Another part of the satellite signals transmit through the waterproof part to be reflected by the dish reflector to the air. However, due to dielectric constants and impedance differences between the waterproof part and the air, there is a reflected wave generated at an incident interface of the waterproof part, which is reflected backward to the feedhorn. In such a situation, a radiating efficiency of the feedhorn is decreased, and a useful bandwidth of the feedhorn may become narrower.\nIn addition, a return loss of the satellite signal or the feedhorn (i.e. a ratio of incident and reflected waves) is relative to the radiating efficiency and the useful bandwidth. Under some conditions, the waterproof part may improve the return loss but narrows the useful bandwidth, and thus the return loss and the useful bandwidth cannot be improved at the same time. Therefore, how to improve both of the return loss and the useful bandwidth of the satellite signals or the feedhorn has become a topic of the industry."}
{"text": "Tubular endoprosthesis or “stents” have been suggested for dilating or otherwise treating stenoses, occlusions, and/or other lesions within a patient's vasculature or other body lumens. For example, a self-expanding stent may be maintained on a catheter in a contracted condition, e.g., by an overlying sheath or other constraint, and delivered into a target location, e.g., a stenosis within a blood vessel or other body lumen. When the stent is positioned at the target location, the constraint may be removed, whereupon the stent may automatically expand to dilate or otherwise line the vessel at the target location.\nAlternatively, a balloon-expandable stent may be carried on a catheter, e.g., crimped or otherwise secured over a balloon, in a contracted condition. When the stent is positioned at the target location, the balloon may be inflated to expand the stent and dilate the vessel.\nBalloon-expanded stents tend to be relatively stiff and straight, as are the balloons used to deliver them, which reduces the ability of the stents to conform to the shape of vessels that are curved and/or angulated. Curved connectors between rings of certain stent designs may allow bending of the unexpanded stent, but such connectors rarely provide enough differential lengthening to allow significant bending of the expanded stent because the connectors are made from the same inelastic material used throughout the stent. Moreover, if such a fully expanded stent were capable of bending easily, e.g., to accommodate a bend in the artery, the stent may be capable of bending repeatedly in response to arterial motion, increasing risk of the stent becoming work-hardened and/or breaking, e.g., after deployment within a body lumen, such as a cardiac vessel, within which the stent may experience significant dynamic forces. The interface between the end of a substantially stiff, straight balloon-expanded stent and a relatively soft, otherwise curved artery may also become the focus of stress. The resulting micro-trauma may cause inflammation, scarring, and/or flow-limiting narrowing, especially when the artery stretches or bends repeatedly with the cardiac cycle, respiratory excursion, and/or flexion/extension of a joint.\nOne solution involves the implantation of many short unconnected stents, so that the stented artery can bend, just as a long train bends. The simultaneous delivery of multiple short balloon-expanded stents is complicated by the tendency of individual stents to migrate relative to the balloon during inflation. When a conventional balloon on a balloon catheter is expanded, one end of the balloon may initially expand before the other, which may cause the stent to migrate away from the initially expanding end and/or compress the stent axially, or both ends may expand initially before a central region carrying the stent, which may cause the stent to compress or otherwise deform undesirably. If this occurs, the actual position of the stent may be difficult to control, which risks the stent being deployed misaligned relative to a desired location. This aspect of balloon expansion may be particularly problematic when deploying many short stents.\nAccordingly, an improved apparatus and methods for delivering stents would be useful."}
{"text": "The present invention relates to an access authorization system comprising at least one arrangement of a plurality of microparticles. In addition to simple mechanical access controls, such as for example keys, combined mechanical and electronic access controls or wholly electronic access controls are also realized in conjunction with access authorization systems, in particular computer systems or computer networks but also access controls for buildings.\nIn this context, queries are made of, for example, access cards or passwords.\nHowever, having proven problematic for electronic access controls is the fact that experts can circumvent such access controls, whereby these experts take advantage of vulnerabilities in the electronic access controls and “hack into” them.\nThere is thus the need to make access authorization systems even more secure, in particular to the effect of preventing or obstructing unauthorized persons from gaining access to protected systems by means of access authorization.\nColor-coded microparticles consisting of a plurality of colored layers in a preselected color sequence, wherein the color sequence represents an identification code, are already known from DE 26 51 158 A1. The microparticles are thereby made of melamine-alkyd resin, whereby, for example, 7 layers of color are deposited one atop the other at a thickness of approximately 100 μm on a polyester carrier film of approximately 50 μm.\nDE 198 53 764 A1 relates to a system for securing and marking products using microparticles which each have a respective plurality of color layers forming a code."}
{"text": "People's demands for energy continuously increase with the progress in various technological fields. To avoid shortage or even depletion of energy resources, endeavors have been made to discover new energy sources. On the other hand, effective means have also been developed to minimize the power consumption of existing electric and electronic devices.\nLighting fixtures account for a very large part in all kinds of power-consuming devices. Presently, the bases for general lamps are mainly divided into two types, namely, conventional and electronic lamp bases. The conventional and the electronic lamp base all are equipped with a starter, but have a conventional and an electronic ballast, respectively. The conventional ballast and the electronic ballast are different in their wiring configuration, fundamental characteristics, use manners and power consumption. The biggest difference between the conventional and the electronic ballast is that the conventional ballast uses a low frequency of 60 Hz, while the electronic ballast uses a high frequency ranged between 30,000 and 50,000 Hz.\nSince the lighting fixtures account for a very large part in the power-consuming devices, improvements of the power consumption efficiency of lighting fixtures and development of creative use manners thereof have become new and important issues. Among others, the use of LEDs to replace the conventional illumination light sources has become a frequently adopted technical means for improvement of various kinds of lighting fixtures. In recent years, the LED-related technique develops quickly to largely widen the applications of LEDs. There are more and more daily necessaries with LEDs. The LED illuminates light using high current with low voltage, and has the advantages of long service life, low power consumption and low heat production. With these advantages, LEDs are environmentally friendly and attract many manufacturers to the research and development of LED tube lights.\nDriving circuits for LED tube lights are generally divided into three types according to their technical features. In the first type of driving circuit for LED tube light, two circuit boards are included, one of which is a power driving circuit board while the other one is an LED circuit board. The power driving circuit board and the LED circuit board are electrically connected to each other via conductors or connection terminals and all are mounted in the tube light, so that power can be effectively supplied from the power driving circuit board to the LED circuit board to produce light sources. In the second type of driving circuit for LED tube light, three circuit boards are included in the LED tube light. Two of the circuit boards are power driving circuit boards while the third one is an LED circuit board. The two power driving circuit boards are separately arranged at two opposite ends of the tube light, and the LED circuit board is arranged in the tube light at a middle section thereof and electrically connected via conductors or connection terminals to the power driving circuit boards located at two opposite ends of the tube light, such that power can be effectively supplied from the power driving circuit boards to the LED circuit board to produce light sources. In the third type of driving circuit for LED tube light, only one LED circuit board is included in the LED tube light. The LED circuit board is connected via conductors to two pinned bases mounted to two opposite ends of the tube light, and a power driving circuit is arranged outside the tube light or in a corresponding lamp base.\nIn either one of the above-mentioned three types of driving circuits for LED tube light, complicated wiring is required to achieve electrical connection between the LED circuit board and the power driving circuit boards or the externally arranged power driving circuit. Since the tube light has only very limited internal space, complicated working procedures are needed to mount the LED circuit board and the somewhat bulky power driving circuit boards in the small internal space of the LED tube light, which inevitably largely increases the manufacturing and material costs of the LED tube light."}
{"text": "Various embodiments of the invention relate to a focus detecting apparatus.\nAs digital photographing apparatuses, such as digital cameras, camcorders, or the like, have been miniaturized and as technology relating to a battery, or the like, has been developed, digital photographing apparatuses can be easily carried. Thus, an image can be easily captured anywhere. In addition, digital photographing apparatuses provide a wide variety of functions that enable non-professionals to capture a good quality image.\nIn order to capture a good quality image of a subject, light has to be sufficiently irradiated onto the subject. When light is not sufficiently irradiated onto a subject, it is not easy to focus on the subject and a captured image is dark, and thus it is not easy to recognize the photographed subject. Thus, a lighting device for lighting a subject may be embedded in a digital photographing apparatus or may be separably installed on the digital photographing apparatus as occasion demands.\nIt is required that a subject be precisely focused on in order to capture a clear still image or a clear moving picture image by using a digital photographing apparatus, such as a camera, a camcorder, or the like. Examples of autofocusing (AF) methods include a contrast AF method and a phase difference AF method.\nIn the contrast AF method, a photographing operation is performed when a position of a focus lens is changed, a contrast value is obtained from an image signal generated by an image capturing sensor, and the focus lens is driven to a position at which the contrast value is a peak.\nIn the phase difference AF method, an image capturing sensor and an additional sensing device are disposed, and a focus position is detected from a phase difference in beams of light irradiated onto the sensing device."}
{"text": "Swimming pool cleaning devices (hereinafter pool cleaners) are used for maintaining residential and commercial swimming pools in a clean and attractive condition. Pool cleaners have been developed for cleaning and/or dislodging settled debris from the floor and side wall surfaces of the swimming pool, thereby substantially reducing the need for manual vacuuming and/or brushing of the floor and side wall surfaces of the swimming pool.\nSome swimming pool cleaning devices may be powered by a floating power supply which floats on a water surface of the swimming pool. A power cable may be used to connect the swimming pool cleaning device to the floating power supply. The floating power supply is generally coupled to a main power supply source located external of the swimming pool. A main power cable may be used to couple the floating power supply to the main power supply source.\nUnfortunately, the main power cable is generally laid across the deck of the swimming pool. Not only is this unsightly, but it poses a hazard to those around the swimming pool. The main power cord may pose a tripping hazard to those walking around the swimming pool. Further, since the main power cord generally has to be long enough to not limit where the swimming pool cleaning device can be used in the swimming pool, those around the swimming pool may get tangled within the main power cord thereby posing not only a tripping hazard but a potential choking hazard.\nTherefore, it would be desirable to provide a system and method that overcomes the above."}
{"text": "Applications have a user interface that can take a variety of different forms. One common user interface includes menus and toolbars. On a computer system, each application program can have its own user interface which appears when the application program is executing. An application may also be a composite application consisting of multiple, different application programs. Each of the different application programs may be a component of the composite application. The composite application may present to a user a single common interface for all the component application programs. With this single interface, all the application programs of the composite application may share the same command menus and toolbars. A single user interface for the composite application may be shared between all the application programs. One drawback of the foregoing is that each application program contributing to the single user interface may be developed independently of the other application programs of the composite application. In such instances, the user interface for each program of the composite application may require a collaborative effort while the remaining portions of each application program are independently developed."}
{"text": "In some printing apparatuses, toner is applied to a substrate to form a toner image. The image can be heated while being subjected to pressure by a fixing device to fix the toner to the substrate. In these apparatuses, the fixing device can be subjected to temperature conditions that shorten the lifetime of components of the fixing device.\nIt would be desirable to provide fixing systems and methods for fixing marking material to a substrate that can utilize temperature conditions that allow lower run costs and desirable image quality."}
{"text": "Generally, if the position of an object acting as a communication object is changed or if scanning is needed to search for the position of the communication object, a directional pattern of an antenna must be changed.\nConventional art controls the direction of a main beam by changing a phase difference between array radiation elements, or changes a directional pattern using mechanical rotation.\nHowever, according to the conventional art in which the phase difference between array radiation elements is changed, a plurality of additional circuits may be needed to control the phase of each array radiation element, the angle of pattern variation is a small angle, and a high side lobe occurs, resulting in a reduction of radiation efficiency of each antenna.\nIn addition, according to the conventional art in which the directional pattern changed using mechanical rotation, a separate structure for rotating the antenna is needed. If the communication object moves at high speed, it may be difficult for the directional pattern to be changed in the accurate direction."}
{"text": "1. Field of the Invention\nThe present technology relates to an information processing apparatus creating drawing data based on a printing job, and the like.\n2. Description of the Related Art\nSo called production printing that is of printing service, in which a lot of documents for business use are printed or bound, is known (for example, Japanese Laid-open Patent Publication No. 2012-238188). In the Japanese Laid-open Patent Publication No. 2012-238188, a printing system for informing a user whether a post process is available or not, in view of the whole system, is disclosed.\nIn the production printing, a printing process is usually handled as a workflow, and opening of a printing workflow is proceeded with. By opening the printing workflow, it is possible to describe setting files, or the like for printing jobs in important printing processes, etc., by common format, in software (workflow applications, described below), or printers of different manufacturers. A standard format referred to as a JDF (Job Definition Format) is known, as a format for describing all the processes of the printing workflow.\nThere are various processes in the printing workflow, such as a process for creating documents or contents, processes for designating printing methods, printing processes, post processes, or the like. Although these processes are performed by various workflow applications or printers, cooperation or printing process management can be achieved between printers by the JDF.\nHowever, some workflow applications or printers may extend the format of the JDF. In this case, the JDF provided by the workflow applications of the different manufacturers may include a description in proprietary format.\nFIG. 1A is an illustration diagram for illustrating an example of an inconvenience in accepting printing jobs by a print processing device of a manufacturer “C”. Additionally, FIG. 1A and FIG. 1B shows an example of comparison, not an example of prior art. The print processing device of “C” receives printing jobs from workflow applications of a manufacturer “A” and a manufacturer “B”. It may occur that the print processing device of “C” cannot analyze the JDF to process since the respective workflow applications extend the format of the JDF.\nTherefore, in order to accept the printing jobs of the respective workflow applications, the print processing device of “C” may have a rendering engines 59 (hereinafter, referred to as RIP engines 59) capable of JDF conversion and handling respective workflow applications. In FIG. 1B, a JDF analyzing unit 56 analyzes the JDF to determine the manufacturer of the workflow application, and converts the setting information so as to be processed by the print processing device of “C”. Also, the printing jobs can be processed to print with user's desired finished appearance, by installing the RIP engines 59 capable of handling respective workflow applications.\nThus, the print processing device of “C” can process the printing jobs to print with user's desired finished appearance, even if the respective workflow applications extend the format of the JDF.\nBy the way, a user may need to display or change contents of the printing jobs before rendering the printing jobs by the RIP engine. However, since settable attributes or settable range of values of the attributes may differ by respective RIP engines, the print processing device of “C” may not appropriately display the content of printing jobs of “A” or “B”. Further, an instruction, to change the value of the attribute into a value which is valid only for the printing job of “B”, may be accepted as an instruction for the printing job of “A”.\nAn example of aggregation printing imposition will be described. In the workflow application of “A”, only one of setting options for aggregation printing imposition of 2-up (2 pages are aggregated into 1 page)/4-up (4 pages are aggregated into 1 page)/9-up (9 pages are aggregated into 1 page)/16-up (16 pages are aggregated into 1 page) is acceptable, and the setting options are displayed as selectable options to accept one of the options. Meanwhile, in the workflow application of “B”, the aggregation printing imposition can be set by using a format of “M (number of pages in longitudinal direction)×N (number of pages in lateral direction)”, where any combinations are accepted as far as “M” and “N” are set within the respective limits.\nIn the workflow application of “C”, the same setting screen as that of the workflow application of “A” is used. If the print processing device of “C” displays the printing job of “B”, in a case where the aggregation printing imposition set as “M×N”=“1×2”, “2×2”, “3×3”, or “4×4” can be displayed respectively as 2-up, 4-up, 9-up, or 16-up. However, in a case where the combination of “M” and “N” is not one of the options shown above, the print processing device of “C” cannot appropriately display the setting information of the aggregation printing imposition.\nTo the contrary, a case, where the same setting screen as the workflow application of “B” is used in the workflow application of “C”, will be described. In this case, the print processing device of “C” can display the printing jobs of 2-up, 4-up, 9-up, or 16-up, and accept an instruction to change into any one of the combinations of “M×N”. However, in a case where the changed combination of “M×N” is neither “1×2” nor “M”=“N”, the rendering cannot be performed since the RIP engine of “A” does not support such an aggregation printing imposition."}
{"text": "The present invention relates to radio frequency (RF) power amplifiers (PA) module. Portable devices such as laptop personal computers (PC), Personal Digital Assistant (PDA) and cellular phones with wireless communication capability are being developed in ever decreasing size for convenience of use. Correspondingly, the electrical components thereof must also decrease in size while still providing effective radio transmission performance. However, the substantially high transmission power associated with RF communication increases the difficulty of miniaturization of the transmission components.\nA major component of the wireless communication device is the radio frequency PA. The PA is conventionally in the form of a semiconductor integrated circuit (IC) chip or die in which signal amplification is effected with substantial power. The amplifier chip is interconnected in a circuit with certain off-chip components such as inductors, capacitors, resistors, and transmission lines used for controlling operation of the amplifier chip and providing impedance matching of the input and output RF signals. The amplifier chip and associated components are typically assembled by interconnected metal circuit and bond wires on a printed circuit board (PCB) having a dielectric substrate or a lead frame.\nAmong significant considerations in the miniaturization of RF amplifier circuits is the required impedance matching for the input and output RF signals of the amplifier. Input and output impedance matching circuits typically include capacitors, resistors, and inductors in associated transmission lines or micro strips for the RF signals into and out of the amplifier chip. However, these impedance matching circuits may require specifically tailored off-chip components located around the amplifier IC chip. Accordingly, the application circuitry must include many electrical input and output terminals or bonding pads to which the corresponding portions of the off-chip impedance matching circuits are separately joined. This increases the difficulty of assembly and required size of the associated components, and affects the overall manufacturability of the portable devices.\nOne important requirement for the state-of-the-art wireless devices is to provide power amplification in a plurality of frequency bands. The quality and power level of the amplified RF signals need to be properly controlled. The amplification of RF signals is required to be linear over a wide signal power range in each of the plurality of frequency bands. Preferably the amplification is reduced or increased according to input RF signal, transmittance range and data rate so that power consumption can be optimized."}
{"text": "This invention relates to a device for transporting and positioning dough triangles in crescent shaped dough rolls forming machines.\nThe technique for preparing crescents, also called \"croissants\", consists in preparing a strip of dough which is rolled on a roller machine.\nAfter some dough strip calibration operations, the strip is transferred into a machine which cuts out triangles.\nIn order to waste no materials and reduce costs, the triangles are arranged, after the cutting thereof, in parallel rows with opposing orientations, as shown in FIG. 1.\nAfter cutting, the triangles or at least one half of them must be orientated such that they are all presented to the rolling machine with their bases onwards.\nOn commercially available machines, these orientation operations are carried out by simply turning upside down alternately one half of the triangles, as shown diagramatically in FIG. 2.\nThe triangles will then enter the rolling machine which comprises essentially a main roller A which carries the dough triangle 2, an upper roller B which guides the dough triangle 1, and two roll-up belts C and D which perform the rolling operation with the aid of the roller A (FIG. 3).\nThe problem encountered with this processing originates from the fact that the dough, upstream of the cutting station, is located on a continuous conveyor belt, thereby the top face, being exposed to air, is drier than the bottom face which bears onto the belt and is thus prevented from losing moisture.\nThis position is also satisfactory on the roll-up machine, because the wetter face will adhere on the roller A which transfers it onto the roll-up belts C and D without problems, since a weak adhesion engagement is established between the dough and roller B.\nHowever, when the triangles 1 which have been upturned arrive, the higher adhesion due to higher moisture will occur on the roller B, so that the dough triangle 1 readily separates from the roller A and is not inserted in between the roll-up belts C and D and is instead ejected, as shown in dotted lines in FIG. 3.\nThis situation produces considerable inconvenience, accompanied by a reduced output, and requires constant attention by an operator for recovering the high number of dough triangles which are not processed."}
{"text": "German patent DE 195 25 131 C2 discloses a known method. In this known method, a sensor having two measuring coils is used, which are arranged immediately one after the other perpendicular to the conducting surface, and also offset from each other. In the known method, frequencies are applied to the input side of the two measuring coils from an oscillator in each case, where the oscillator frequencies differ from each other. Depending on the respective distance of the sensor from the conducting profiled surface, there is a resultant impedance change in the measuring coils that is detected by means of differential amplifiers connected to the output side of the measuring coils.\nA demodulator unit comprising two demodulators and connected to the output of the differential amplifiers, and a circuit comprising two EPROMs connected to the output of the demodulator unit, are used to generate digital measured-values that correspond to the distance of the respective coil of the sensor from the conducting profiled surface. The digital measured-values are supplied to an arithmetic unit implemented in the form of a comparator and then examined there to evaluate whether the digital measured-values of the two measuring coils agree within a certain tolerance range. If this is the case, then the distance currently detected is deemed to be measured correctly, and an OK signal is generated by the arithmetic unit."}
{"text": "Vehicles, such as automobiles and trucks for example, include a rear drive module (RDM) that is connected to the vehicle engine by a prop-shaft. The prop-shaft transmits rotational energy (torque) developed by the vehicle engine to the rear drive module, which in turn transmits the rotational energy to the wheels. In a rear-wheel drive vehicle, the prop-shaft directly couples the RDM to the vehicle's transmission. In an all-wheel or four-wheel drive vehicle, additional components may also be included, such as a power take-off unit for example.\nIt should be appreciated that the transmission of rotational energy from the propshaft to the RDM, and from the RDM to the wheels generates reaction forces within the RDM to counter the transmitted torque. These reaction forces generally may be characterizes as a “roll” type and a “pitch” type movement. The roll type movement is a rotation about a longitudinal axis passing through the RDM. A roll type movement may cause the axles to flex with respect to the RDM and cause undesired noise and vibrations. A pitch type of movement is a rotation about the lateral axis the RDM due to a reaction to the drive torque at the wheels. Articulation of the RDM due to pitch also results in undesirable noise and vibration, and may also reduce the operating life of the prop-shaft.\nTraditionally, to counter the reaction torques placed on the RDM, a mounting system was used that securely coupled the RDM to the vehicle structure, such as directly to the vehicle frame, or to an intermediary cross-member or cradle-member. Typically, these systems used some type of three-point mount that included isolation bushings that reduced the transmission of vibration from the RDM to the structure. It should be appreciated that these vibrations may have been due to the operation of the RDM and by the operation of the engine as well.\nTraditionally, the vehicle engine was an internal combustion engine having cylinders that are alternately fired to produce the rotational energy. Due to a need to improve fuel efficiency, alternate control schemes for the vehicle engine have been developed that selectively deactivate cylinders. Under certain circumstances, when a cylinder is deactivated, no fuel is combusted and fuel efficiency is increased. However, it has been found that the deactivation of cylinders results in low frequency vibrations being transmitted to the RDM via the prop-shaft that were not previous experienced in traditional engine control configurations. Further, it has been found that in some circumstances, existing RDM mounting arrangements were inadequate to counter the excitation forces generated at these low frequencies.\nAccordingly, it is desirable to provide an RDM and RDM mounting arrangement that provides a desired level of performance when subjected to low frequency vibrations from the vehicle engine."}
{"text": "The present invention relates generally to landfills, and more particularly to systems and methods for disposing of liquid condensate from landfill gas recovery systems.\nWaste products decompose in landfills, and after the free oxygen in the landfill is depleted, the waste product decomposition generates methane gas. It is desirable to recover this methane gas for environmental and safety reasons. To this end, landfill gas recovery systems have been introduced which collect the gas generated in landfills and burn the gas in flares on the landfill.\nOccasionally, gas in the recovery system condenses with other fluids such as water. This methane-based condensate, like the gas, must be removed from the landfill for safety and environmental reasons, and to ensure that blockage of gas piping and damage to the flare system does not occur. Typically, the condensate is simply pumped out of the gas recovery system and transported to a hazardous waste dump site, where it is disposed of.\nAs recognized herein, transporting hazardous condensate to another waste facility for disposal is not only expensive, it does not solve the environmental problem of disposing of the condensate, but rather only moves the problem to a hazardous waste disposal facility. With this in mind, the present invention recognizes the desirability of economically disposing of the condensate at the site at which it is recovered in an environmentally benign way.\nAs recognized herein, one method for disposing of the condensate is to burn it in the flare chamber that is used to burn the methane gas. Typically, a landfill gas recovery flare chamber includes a ring of vertically-oriented burners located near the bottom of the chamber, and methane gas is piped through the burners and oxidized, with the hot oxidation products exhausting upwardly up through the flare chamber and out of the open top end of the chamber. In such a flare chamber, the condensate can be injected radially into the flare chamber above the burners by entraining the condensate in a pressurized high velocity air stream above the flame of the flare.\nSuch a system, as understood by the present invention, unfortunately requires a relatively expensive air compressor to generate the pressurized air stream. Also, a portion of the high velocity condensate stream tends to impinge on the wall of the flare chamber that is opposite the condensate injection point, damaging the wall.\nAlternatively, the present invention understands that condensate can be pumped upwardly into the flare chamber through a vertical pipe that is centrally located in the flare chamber below the ring of burners. As the condensate moves upwardly past the burners, it flashes into vapor. As recognized by the present invention, however, the injection rate of condensate sometimes must undesirably be limited to avoid excessively cooling the flare chamber as the latent heat of vaporization of the condensate is overcome. Excessively cooling the flare chamber could reduce the ability of the flare to burn the methane gas and condensate. Moreover, the present invention understands that landfill process controls, including those related to condensate injection systems, preferably be automatic, to more accurately control the processes and to avoid the necessity of personnel undertaking time consuming and repetitive process monitoring and adjustment.\nAs further recognized herein, it is possible to provide a condensate injection system having a relatively high condensate injection rate without excessively cooling a flare chamber, and to automatically control the condensate injection rate as appropriate for the particular energy level of the flare. Accordingly, it is an object of the present invention to address one or more of the abovenoted considerations.\nA compressorless condensate injection system is disclosed for a landfill having a flare chamber including at least one wall that is heated when the flare chamber burns methane gas extracted from the well. The system includes a condensate reservoir and a condensate pump in fluid communication with the reservoir to pump condensate into the chamber at a high pressure, preferably 40-250 pounds or more. At least a first injection line is in fluid communication with the condensate pump but not with an air compressor. The first line terminates in a first nozzle that is positioned on the flare chamber for directing condensate into the chamber such that condensate from the nozzle is vaporized when it is sprayed into the chamber without requiring the use of compressed air.\nIn a preferred embodiment, the first line has a heat exchange segment that is curved, e.g., the segment can extend partially or completely around the flare chamber before terminating in a nozzle. In this way, fluid in the first line can be heated when the flare chamber burns gas extracted from the well.\nA first control valve preferably is in fluid communication with the first injection line for selectively blocking fluid flow therethrough, with the first control valve being responsive to electrical control signals. Indeed, secondary injection lines with respective solenoid valves and nozzles can be provided for selectively injecting even greater amounts of condensate into the chamber, depending on vaporization conditions. These secondary nozzles can be oriented to direct condensate upwardly and radially inwardly into the flare chamber. If desired, a ring line can communicate with the condensate pump, and the ring line terminates in a ring line nozzle disposable adjacent the burners of the flare.\nAdditional features can include a methane gas inlet line and a methane sensor for measuring a methane concentration in the inlet line, a flow sensor for measuring gas flow rate in the inlet line, and a temperature sensor for sensing temperature in the flare chamber. Also, condensate temperature and pressure can be measured in each heat exchange segment. Electrical control signals for controlling the solenoid valves can be generated by a computer based on these signals.\nIn another aspect, a computer program device can include a computer program storage device readable by a digital processing system, and a computer program on the program storage device and including instructions executable by the digital processing system for performing method steps for controlling at least one control valve disposed in at least one condensate injection line in a landfill flare chamber. The method undertaken by the computer includes determining a gas volume burn rate based on a combination of methane concentration in gas to be burned in the chamber, flow rate of gas, and flare chamber temperature. Also, the computer generates one or more control signals to control the valve or valves in response to the determination of gas volume burn rate.\nIn still another aspect, a condensate injection nozzle includes a nozzle body defining a pathway therethrough, and an orifice element disposed in the pathway. An diversion plate is also disposed in the pathway. In accordance with present principles, the diversion plate causes turbulent flow of the condensate, prior to the condensate passing through the orifice element and being injected into the flare chamber.\nThe details of the present invention, both as to its structure and its operation, can best be appreciated in reference to the accompanying drawings, in which like reference numerals refer to like parts, and in which:"}
{"text": "1. Field of the Invention\nThe invention generally pertains to the field of remotely located network connected intelligent devices.\n2. Description of the Prior Art and Related Information\nIn legacy bandwidth-limited distributed networks prior to the Internet era, to update the code of a large numbers of (lottery, for example) terminals, a download server typically “pushed” the data to each terminal. Under such a scheme, scheduling and error recovery are carried out entirely under the control of the download server.\nHowever, new generation lottery terminals, gaming machines and Point-of-Sale (POS) terminals based on PC architecture or other multimedia-enabled architecture may require frequent and voluminous updates and downloads of programs and data in order to provide continuously updated rich services. In such systems, downloading is commonly carried out using a traditional “pull” method in which each remote machine is scheduled to initiate a download at a predetermined time from a predetermined remotely located server. At the scheduled time, the entire transfer is carried out under the control of the remote terminal, including error recovery. Well-known and popular downloading utilities include programs such as GetRight (www.getright.com) and Gozilla (www.gozilla.com). Using such programs, however, the server that delivers the data file to be downloaded by the remote terminal devices is usually a generic FTP server that does not have capability of intelligently managing the network traffic.\nGeographically distributed download cache technologies, such as available from Akamai (www.akamai.com) and Digital Island “2Deliver” service (www.digitalisland.com) accelerate Internet network performance when downloading static data from global Internet servers. For example, Amakai has deployed on the order of 10,000 servers around the world.\nThere is a significant risk for the network that links the remote terminal devices and the central system to be subjected to unauthorized intrusion, virus infection and distributed denial of service (DDOS); consequently costly bandwidth limited private networks are often preferred. Alternatively, Virtual Private Networks (VPN) to carry secure communication through an encrypted tunnel via the Internet is becoming increasingly popular for company inter-communications. However, the setup and infrastructure management costs are high, as is the cost of training software developers. Moreover, the scalability of VPNs to very a large number of client devices has not been demonstrated at this time.\nUncontrolled data downloads may render the operational network traffic useless, with the same consequences as Distributed Denial Of Service attacks (DDOS); therefore, data downloading is conventionally scheduled outside of operational hours. The requirement to perform data downloads outside operational hours results in significant waste of data bandwidth resources and longer download campaigns.\nIn addition, with the traditional “pull” download method, there is no feed-back that would enable performing a close-loop regulation of the individual terminal device download rates in order to ensure a uniform or predetermined download level. Although the “push” method allows fine-grained download bandwidth throttling, error recovery management requires a very complex download server that cannot easily scale to a very large number of client devices. Legacy distributed network also make use of broadcast download techniques, but error recovery is complex and the Internet infrastructure cannot readily support such broadcast download techniques.\nFIG. 6 is a flowchart illustrating a conventional unregulated download session. In FIG. 6, the boundaries between the remote device (such as devices 104-122, for example), the network 124 and tire transactional server 102 of FIG. 1 are shown in dashed lines. According to the conventional method of FIG. 6, a device whose identifier (ID) is XYZ initiates a download session to retrieve data file ABC from a server 102, at step S601. The download session then requests the first packet of file ABC, as shown at S602. The server 102 receives this request, opens a download session for file ABC at S603, retrieves packet #1 and sends the retrieved packet #1 to the requesting remote device at S605. The remote device receives packet #1 from the server and stores it, as shown at S606. Packet #2 is then requested, retrieved and sent back to and stored by the remote device, as shown at S607, S608, S609 and S610. The remote device then continues in a similar manner, until Packet n is requested at S611 and retrieved at S612. As the transactional determines that Packet n is the last packet of file ABC, the server 102 sets a last packet flag at S613 and sends the last Packet n together with the last packet flag at S614. The remote device then receives this last Packet and stores it, as shown at S615. Upon receipt of the last packet flag together with the last Packet n, the remote device closes the download session, as shown at S616. Upon confirmation of the good receipt of file ABC by the remote device, the server 102 may then close the download session opened in step S603.\nAs shown in FIG. 6, the remote device initiates the download session, which session is ended by the remote device when the last data packet is received. The remote device, therefore, is called the session master. Details of error recovery are not shown in FIG. 6. However, it will be apparent to those of skill that the remote device may request the transmission or re-transmission of any packet. The flow of data in FIG. 6 is un-regulated and download speed is inherently limited by the capacity and congestion of the network 124. Each remote device “fights” to get its own data immediately. When considering a very large number of download sessions using the same unregulated scheme, data traffic can be drawn down to a crawl, thereby denying high priority traffic and resulting in the dreaded DDOS (Distributed-Denial-Of-Service). Although QOS (Quality-Of-Service) mechanisms are available in order to route traffic according to priority attributes, such mechanisms are not universally implemented in routers across all wide area networks. Consequently, prioritization of traffic using QOS or other similar scheme by routers does not regulate data traffic in a satisfactory manner."}
{"text": "The proliferation of electronic devices capable of displaying video information in the form of movies, television programs and games, for example, has prompted great demand for video content. The advent of low cost video recording devices and the ability to circumvent anti-copy protection techniques has lead to unauthorized copying and distribution of such video content. In an effort to reduce the incidence of unauthorized copying, content creators now place a forensic mark, sometimes referred to as a watermark, within the video content for identification. By tracing the forensic mark, the content creator can isolate the source of the unauthorized copies.\nThe forensic mark can take various forms. For example, a content creator can apply a unique serial number to each copy or to a batch of copies. Alternatively, the content creator could apply a unique combination of alphanumeric characters or graphical symbols. The nature of the mark will depend on various factors, including but not limited to, the nature of the content itself.\nPresent day techniques for forensically marking video content suffer from several difficulties. Placing a forensic mark within the image itself incurs the disadvantage that the forensic mark becomes easy to spot and easy to circumvent by either editing or applying pixilation to the mark for example. Altering the video format to create a forensic mark incurs the disadvantage that converting the video content from one format to another can destroy the mark.\nThus, a need exists for a technique for forensically marking video content that overcomes the aforementioned disadvantages."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrical connector assembly, especially to an electrical connector assembly having electrical connector with low profile and processor having cone pins to make it be assembled to the electrical connector with little force.\n2. Description of the Related Art\nU.S. Pat. No. 6,135,784 issued to Pei on Oct. 24, 2000 discloses a LIF Pin Grid Array (PGA) electrical connector. The electrical connector comprises an insulative housing with a plurality of contacts received therein and a processor assembled on the insulative housing. The processor includes a plurality of pins connecting with the contacts. The contact comprises a pair of spring arms. When the processor is assembled to the insulative housing, the pin of the processor exerts a force on the two spring arms pushing the spring arms move away from each other, and finally the pin connect with the spring arms. Thus, an electrical connection is established between the processor and the electrical connector. However, the number of the contacts becomes more and more, and the force that exerts on the processor becomes larger and larger. Thus, it is difficult to assemble the processor and the pin is easily deformed.\nU.S. Pat. No. 6,544,064 issued to McHugh on Apr. 8, 2003 discloses a ZIF Pin Grid Array (PGA) electrical connector. The electrical connector comprises a base with a plurality of contacts received therein, a cover assembled to the base and a lever assembled between the cover and the base. When the processor is assembled to the electrical connector, the processor and the cover are in the first position that the pins of the processor disconnect with the contacts. Then rotate the lever to make the cover move relative to the base and the processor also moves with the cover to a second position that the pins connect with the contacts. When assembled the processor to the electrical connector, this type of electrical connector is easy to operate and can prevent the deformation of the pins. However, the cover and the lever occupy more space, which violates the development trend of the minimization of the electrical connector.\nHence, an improved electrical connector assembly is required to overcome the disadvantages of the prior art."}
{"text": "Generally, a DRAM (Dynamic Random Access Memory) device includes a plurality of unit cells, each of which comprises a MOS transistor and a storage capacitor. In continuing the trend of higher memory capacity, the size of the unit cell has been continuously decreased in order to increase the packing density of the DRAM device. The reduced cell size results in a decrease in capacitor area of the unit cell. The decreased capacitor area means low cell capacitance, which often induces problems like as low read-out capability and soft error.\nOne proposal to solve the above-mentioned problems is to use high dielectric constant material as a capacitor dielectric layer, which constitute a capacitor together with a lower electrode and an upper electrode. Typical examples of the high dielectric constant material are tantalum oxide (Ta2O5) and BST ((Ba,Sr)TiO3). The material of the lower electrode or the upper electrode is required to have a high work function value and not to be reactive with the capacitor dielectric layer. A typical example of the material of the lower and the upper electrodes is a noble metal, which includes platinum, ruthenium, iridium, rhodium and osmium.\nFIG. 1 is cross-sectional views illustrating a conventional method for forming a capacitor using tantalum oxide layer as a capacitor dialectic layer. Referring to FIG. 1, a lower electrode 20 is formed on a substrate 10. The material of the lower electrode 20 is ruthenium. A capacitor dielectric layer 25 is deposited on the lower electrode 20. The material of the dielectric layer 25 is tantalum oxide. The thickness of the dielectric layer 25 is 140˜160 Å. The dielectric layer 25 is crystallized by a thermal treatment at 700° C. or more. Subsequently, an upper electrode 30 is formed on the crystallized dielectric layer 25, thereby completing a capacitor 30 on the substrate 10. The material of the upper electrode 30 is ruthenium.\nFIG. 2 is a graph showing a change in equivalent oxide thickness value by the crystallization process. The equivalent oxide thickness value represents an effective thickness of a capacitor dielectric layer of a capacitor on the assumption that the capacitor dielectric layer was made of silicon oxide. Therefore, in general, higher equivalent oxide thickness value means lower capacitance per unit capacitor area. Referring to FIG. 2, the vertical axis represents equivalent oxide thickness value. On the horizontal axis of the graph, the reference “NO” means that the crystallization is not performed, and the reference “700° C.” means that the crystallization is performed at 700° C. As shown in the graph, the equivalent oxide thickness value is favorably decreased by performing the crystallization process.\nHowever, the crystallization process has also problems. That is to say, the crystallization process is usually performed at relatively high temperature, thereby generating a lot of grain boundary in the tantalum oxide layer. The grain boundary often acts as a path of leakage current and may induce unfavorable leakage current in the capacitor. Moreover, the crystallization process under relatively high temperature may induce unfavorable deformation of the lower electrode and damage on the tantalum oxide layer.\nMeanwhile, it is thought to be difficult to decrease the equivalent oxide thickness value into 10 Å or less in the conventional method, because both nuclear generation and crystal growth occur simultaneously during the crystallization process.\nAccordingly, the need remains for a method for forming a capacitor having low leakage current as well as low equivalent oxide thickness."}
{"text": "Commercial technologies developed for vehicular “runflat” or “airless” tires include Michelin's Tweels and Resilient Technologies' “non-pneumatic tire.” Both of these use a honeycomb configuration for consumer applications. Runflat technology is also important to military applications. A representative military vehicle may impose 10,000 lbs of vehicle load per tire. To enable the vehicle to operate after tires are perforated by terrain or gunfire damage, passive runflat systems are employed inside the tire. The runflats currently found on the military vehicle are designed to provide mobility for a short time after a tire goes flat, but weigh approximately 100 lbs per tire. In addition to runflat technology, the vehicle utilizes a central tire inflation system (CTIS) to operate effectively across different terrain conditions.\nPoisson's ratio (v), named after Simeon Poisson, is the ratio of the relative contraction strain, or transverse strain (normal to the applied load), divided by the relative extension strain, or axial strain (in the direction of the applied load). Some materials, called auxetic materials, have a negative Poisson's ratio (NPR). If such materials are stretched (or compressed) in one direction, they become thicker (or thinner) in perpendicular directions.\nNPR materials have attracted significant interest due to their unique behaviors. Unlike conventional materials, a NPR material may shrink when compressed along a perpendicular direction. One result of this behavior is that the material can concentrate itself under the compressive load to better resist the load. Thus, a NPR material becomes stiffer and stronger as the amplitude of the load increases. It has also been found that NPR can improve material/structural properties, including enhanced thermal/shock resistance, fracture toughness, indentation resistance and shear modulus. [1-3].\nAuxetic and NPR structures have been used in a variety of applications. According to U.S. Pat. No. 7,160,621, an automotive energy absorber comprises a plurality of auxetic structures wherein the auxetic structures are of size greater than about 1 mm. The article also comprises at least one cell boundary that is structurally coupled to the auxetic structures. The cell boundary is configured to resist a deformation of the auxetic structures.\nThe vast majority of auxetic structures are polymer foams. U.S. Pat. No. 4,668,557, for example, discloses an open cell foam structure that has a negative Poisson's ratio. The structure can be created by triaxially compressing a conventional open-cell foam material and heating the compressed structure beyond the softening point to produce a permanent deformation in the structure of the material. The structure thus produced has cells whose ribs protrude into the cell resulting in unique properties for materials of this type.\nCommonly assigned U.S. Pat. No. 7,910,193, the entire content of which is incorporated herein by reference, describes two- and three-dimensional NPR structures/materials and applications. The negative Poisson ratio effect causes the surrounding material to concentrate into the local area of loading. Consequently, the material becomes stiffer and stronger in the area of the applied load. Moreover, this stiffening behavior is retained under nonlinear, large deformation response.\nNPR structures can react differently under different applied loads. Three unique features of NPR systems include: a) material concentration, b) bulging effect, and c) impact force mitigation. FIG. 1 of the '193 patent illustrates a reactive shrinking mechanism of a NPR material. The unique property of this structure, which includes a plurality of “nested-V” shapes, is that it will shrink in two directions if compressed in one direction. When the structure is under a compressive load on the top of the structure, more material is gathered together under the load so that the structure becomes stiffer and stronger in the local area to resist against the load.\nCommonly assigned U.S. Pat. No. 8,544,515, also incorporated herein by reference, is broadly directed to the use of NPR materials to produce runflat or airless tires. Such materials should provide improved stiffness/weight and survivability, compared to honeycomb, foam, or other cellular materials. Further, because the stiffening behavior of certain NPR structures is retained under nonlinear, large deformation responses and it can be functionally tailored, a runflat tire system based upon such materials should provide performance responses similar to pneumatic tire.\nNPR tires can be tailored and functionally-designed to optimally meet runflat requirements for both military and commercial vehicles. NPR runflat tires may be fabricated using standard materials and simple manufacturing processes, resulting in low-cost and high-volume production. NPR runflat tire designs are fully compatible with Central Tire Inflation Systems (CTIS), while providing a performance equivalent to current military vehicle solutions but at half the weight."}
{"text": "A prior art method for coding phonetic features was formally set up during the U.S. census in 1858. A certain abstraction of phonetic features was considered to be required, in particular because at that point of time, a fairly high fraction of immigrants in first or second generation lived there who named their children with some assimilated names, i.e., first names or last names. For example, the name Schmidt was assimilated to Smith, or Johannson to Johnson. As it was foreseeable that the original way for writing these names would be lost anyhow, the evaluation of the census was stored supplementary with those phonetical features without using any electronic data processing.\nThis method is known by the name SOUNDEX and is applied until now in many technical fields, for example, in reservation systems in which lists of names are searched according to phonetical features only since customer contact is often by telephone call and the typist may not know the proper spelling of the customer's name. The basic algorithm is the same as it was in 1858 and comprises the rules and steps as follows:\n1. Every code according to SOUNDEX contains four characters: a letter followed by three digits whereby in case of missing digits the code is filled up with the digit ‘0’.\n2. The first letter of the character string which is transformed by the SOUNDEX code, for example a name, is taken without any change as the first code element into the SOUNDEX code.\n3. The letters a, e, i, o, u, y, w, h are not coded. All remaining characters are coded as follows:                b, f, p, v are coded to 1        c, g, j, k, q, s, q, z are coded as 2        d, t are coded as 3        l is coded as 4        m, n is coded as 5        r is coded as 6.        \n4. Two subsequent occurrences, i.e., instances, of the same SOUNDEX code is avoided by only taking the first of them.\nBased on the above mentioned rules the following examples result:\nEulerE460GaussG200HilbertH416HeilbronnH416SchmidtS530SmitS530\nA more detailed description of the prior art SOUNDEX coding method is provided by D. Knuth in ‘The Art of Computer Science, Vol. III’, Addison-Wesley, 2nd ed., S. 391–392.\nThe disadvantages of said prior art SOUNDEX coding method can be summarized as follows:\n1. The SOUNDEX code is limited to 1+3=4 characters. As a typically character string has a bigger length, e.g., a length between 6 and 8 characters not all characters can be coded. Thus, as described above the character strings ‘Hilbert’ and ‘Heilbronn’ have the same SOUNDEX code although they sound quiet different, they have a different length and they can thus be described as ‘not related with each other’. In this respect the SOUNDEX method is too ambiguous.\n2. The rules of the SOUNDEX method are not able to transfer longer character groups into their phonetical equivalence. The SOUNDEX method is limited to avoiding the vowels and to the comprehension of consonants having the same phonetic code. The SOUNDEX code is not adapted to particular letter combinations as are for example the English ‘ight’ as in ‘fight’ which has the same phonetic features as ‘ite’.\n3. Today, phonetic features are often stored together with the original data in a data base. Accordingly, the SOUNDEX code properties have to be checked for the persistency and comparability. Such a check yields that the SOUNDEX code has to be stored as a character string as it comprises at least one character and a character compare procedure has to be started against the code whenever the SOUNDEX code is applied in a database search. Compared to a bit string compare procedure this represents a significant loss of performance."}
{"text": "The present invention relates to a method of preparing 2-sulfochloride benzoates, and, more especially, relates to a method of preparing such compounds which comprises diazotizing an anthranilic acid ester to form a 2-diazonium chloride benzoate and then reacting the 2-diazonium chloride benzoate with sulfur dioxide to form the 2-sulfochloride benzoate. This invention also relates to a method of synthesizing saccharin which comprises pyrolyzing a 2-sulfochloride benzoate at elevated temperatures and then ammonolyzing the resulting o-sulfobenzoic anhydride to saccharin.\nSaccharin is variously known as o-benzosulfimide; gluside; benzoylsulfonic imide; and is the anhydride of o-sulfimide benzoic acid having the formula: ##STR1## Saccharin is employed in the manufacture of syrups, medicine (substitute for sugar), soft drinks, foods, and the like, and is a nonnutritive sweetener which can readily be converted to sodium or soluble saccharin, and is a white, crystalline powder. It has an exceedingly sweet taste (500 times that of cane sugar), a melting point of about 226.degree. C. to 230.degree. C., and is soluble in amyl acetate, ethyl acetate, benzene and alcohol; slightly soluble in water, chloroform and ether.\nSaccharin has been made from toluene by the following series of reactions: ##STR2## The imide is converted to the sodium salt to increase the solubility in water. Saccharin can also be prepared by converting a mixture of toluene sulfonic acids into the sodium salt, then by distillation with phosphorus trichloride and chlorine to obtain the ortho-toluene sulfonyl chloride, which by means of ammonia is converted into ortho-toluenesulfamide. This is oxidized with permanganate, treated with acid and saccharin crystallized out. It is reported that the slight bitter taste associated with the saccharin prepared by either of the above methods is caused by the presence of o-toluamide. Moreover, the disposal of the p-toluenesulfonyl chloride obtain by-product in the above processes has also been a problem.\nBoth of these objections have been attempted to be overcome by two more recent processes, the first [A] commencing with thianaphthene (prepared from styrene and sulfur) and the other [B] with anthranilic acid, as follows [see Noller, Chemistry of Organic Compounds, 2nd Edition, pp. 556-7 (1957)]: ##STR3## Nevertheless, it too is reported that even when utilizing these alternate routes there still is a slight bitter taste associated with the saccharin, here probably caused by the presence of certain trace contaminants."}
{"text": "1. Field of the Invention\nThe present invention relates to a microscope apparatus.\n2. Description of Related Art\nRegarding a super-resolution microscope in the related art, STORM (Stochastic Optical Reconstruction Microscopy) is known (for example, refer to United States Patent Application, Publication No. 2008/0032414 and United States Patent Application, Publication No. 2008/0182336). In this microscope, a fluorescent material or a sample having a fluorescent material adhering thereto is used as an observed sample. This fluorescent material has a characteristic of being activated when irradiated with activation light having a predetermined wavelength and fluorescing to be inactivated when irradiated with excitation light having a different wavelength from the activation light later. A fluorescent material is activated with low density by irradiating an observed sample with weak activation light. By emitting excitation light later in order to make only the activated fluorescent material emit light, a fluorescent image is acquired. In the fluorescent image acquired in this manner, images of the fluorescent material that emits light with low density are separated individually. Accordingly, it is possible to obtain the position of the center of gravity of each image. By repeating such a step of obtaining the position of the fluorescent material multiple times, for example, hundreds of times to tens of thousands of times or more to create one image in which all images of the fluorescent material are disposed at their positions, it is possible to acquire a high-resolution observation image.\nIn addition, as a super-resolution microscope apparatus, a STORM based on a three-dimensional method capable of improving not only the resolution within the sample surface (XY plane) but also the resolution in the thickness direction (Z direction) of the sample is known (for example, refer to Bo Huang, et. al. Science 319, 810-813 (2008)). In such a STORM based on the three-dimensional method, the image of a fluorescent material can be made to have an elliptical shape by giving a predetermined astigmatic difference to the image of the sample by inserting a cylindrical lens in the imaging optical system, and the coordinate in the Z direction can be obtained from the ellipticity."}
{"text": "The present invention relates to the textile arts and, in particular, to an apparatus adapted to provide for the removal of lint, dust and other unwanted materials from a circular knitting machine.\nIt is well known and recognized that the processing of textile fibers, including the knitting thereof into fabric, generates large quantities of fiber lint and other debris. This debris often permeates the environment in which the processing equipment is located, and settles on the exposed surfaces of the equipment. It is well recognized that the continued accumulation of the debris, especially on the active machinery elements, can result in unsatisfactory operation, and possible ultimate machinery failure. In addition, the debris can be trapped in the produced fabric, creating defects therein.\nDebris accumulation and contamination is a significant problem in circular knitting machines, which because of their size can generate a large amount of lint and other debris. Such large machines can have a diameter of many feet, utilizing a large number of reciprocating needles and associated devices, including yarn feeds and linkages needed to synchronize the knitting process. Associated with such operation is the generation of an often prodigious amount of lint and other debris.\nIn order to limit the amount of such debris accumulating on the equipment, and to remove the debris therefrom, a variety of fan-type apparatus have been developed. The device set forth in U.S. Pat. No. 5,195,337 to Alan Gutschmidt presents a typical apparatus, As disclosed therein, a cleaning device for a circular knitting machine includes a centrally mounted arm journaled for 360 degree rotation about a central mounting, typically positioned at the center of the knitting machine, A fan is mounted at the distal end of the arm, and oscillates in a vertical plane over a limited angle to direct a flow of air over portions of the knitting machine as the arm rotates."}
{"text": "During the mounting of wheels, the alignment between the hub bore and the bolt holes can rotate making the use of nuts and sleeved cap nuts difficult to align within the bolt holes. The operator must rotate the wheel assembly to align bolt holes while maintaining its position on the center hub bolts. A brake drum is mounted to a wheel hub on the wheel hub mounting tabs, which are typically very short. As the wheels are assembled over the brake drum, the assembly process can cause the brake drum to slide from the wheel hub mounting tabs causing a non-flush connection to the wheel hub itself When the brake drum is non-flush with the wheel hub, the nuts used to attach the wheels to the wheel hub may not be sufficiently threaded onto the wheel studs to securely attach the wheels to the wheel hub and prevent the wheels from falling off the vehicle during operation."}
{"text": "This invention relates generally to semi-automated dialer system and method to provide fast and efficient dialing from a mixed entry sequence. More particularly, the present invention relates to a method for determining intended entry from a mixed dialing sequence that includes both numeric and non-numeric input.\nIn mobile communications appliances equipped with a deterministic input device, such as a QWERTY keyboard, the large number of keys may be a disadvantage. To keep the number of physical keys to a smaller number, it is common to find that the numeric keys share the same physical keys as some of the text symbol keys. For example, in one such implementation on a Windows Smartphone®, the E, R, T and Y are also labeled 0, 1, 2 and 3, respectively.\nThe particular allocation of numerical keys on the alphabetic keypad is flexible by design, but most implementers attempt to construct a layout that approximates the numeric layout typical of a telephone keypad.\nThis ambiguous structure results in design choices that are contextually driven and it is usual to find that the keyboard operates either in numeric mode or in text symbol mode. However, there is limited interactivity between the two modes. Usually it is left to the user to determine which mode is desired.\nCurrent dialer applications which signals a numeric string to the network system via the appliance's transceiver is ill conditioned to send strings that are not entirely numbers.\nSimilarly the fact that the text symbols normally allocated to the touch-tone keypad have no resemblance to the way that numbers are allocated to the qwerty layout results in an unexpected complexity for the typical user. For example, in voice-mail systems it is usual for a caller to be asked to identify the recipient by entering a numeric string that ambiguously spells their surname or first part thereof. If there is an ambiguous possibility where more than one recipient could be intended, the system may resolve this interactively with the user. It is extraordinarily difficult for a user of a qwerty keyboard labeled with a group of numbers, having no other labeling to show the possible ambiguous meaning of the numeric keys in the context of a telephony application, to perform this entry task accurately. This is further exacerbated by the fact that the staggered key layout of a QWERTY keyboard lends a distortion to the numeric labeled keys, as if the standard telephony pad were not rectangular but more a parallelogram shape, which further complicates the ergonomic task.\nIt is therefore apparent that an urgent need exists for an improved system and method for semi-automated dialing using mixed sequences for input that is both accurate and efficient. This solution would fulfill a long felt, yet unmet, need for dialing applications that is able to effectively provide dialing when provided input that contains both numeric and non numeric input; thereby increasing effectiveness of dialing and number entry on a mobile device."}
{"text": "Present day techniques for the spray application of coatings are generally categorized in one of three basic spray application techniques. The first of these techniques is air atomization wherein the coating particles ride on an air stream from the spray gun to the product being coated. The second of these techniques is airless atomization wherein the coating material is atomized and propelled by hydraulic pressure. The third is electrostatic spraying wherein the coating material is atomized by air or airless techniques and the deposition of coating on the product being sprayed is by electrical attraction of the coating particles. With each of these spray techniques spray booths are commonly employed and usually required by federal or state regulatory agencies.\nPresent day spray booths are designed with the objectives of providing a safe working place, preventing pollution of the atmosphere and enhancing the quality of the product being sprayed. The rapid and thorough removal of volatile solvents and vehicle fumes from the premises is essential in many spray applications to meet the requirements of the Occupational Safety and Health Administration and insurance carriers. The air moving through the spray booth carries coating over-spray away from the product, avoiding the finish marring consequences of semi-dry coating particles settling on already coated surfaces. By effectively removing coating particles from the air being discharged to the outside, the spray booth eliminates a common cause of air pollution and thus helps the user avoid risking violations of legal requirements, particularly those of the Environmental Protection Agency, and prevents staining and dirtying of immediately adjacent property.\nNumerous filter designs have been proposed for spray booth applications, but in each case these designs have been found to be inadequate due to the fact that they have either clogged too rapidly or have been cumbersome and costly in design. There exists a need for a replaceable filter for use with spray booths that is simple in construction and suitable for allowing large volumes of air to pass through with a minimal loss of efficiency.\nU.S. Pat. No. 2,909,237 discloses a collapsible filter for separating particulate matter from a fluid stream. U.S. Pat. No. 3,075,337 discloses a replaceable filter medium that is particularly suited for use in spray booths. The filters disclosed in each of these patents, however, are complex in design and construction and do not offer the advantages of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a mechanism for moving an optical pickup radially along a disc. Particularly, the invention is concerned with an optical pickup moving mechanism suitable for the reduction of cost and weight.\n2. Description of the Prior Art\nIn a disc unit which, with use of an optical pickup, records and reproduces information to and from a disc such as CD (compact disc), MD (minidisc), or DVD (digital versatile disc), an optical pickup moving mechanism as shown in FIG. 7 is provided for moving the optical pickup radially along the disc.\nIn the same figure, an optical pickup P0 is principally composed of an objective lens 2, a drive mechanism (not shown) for actuating the objective lens 2, and a base 1 which carries thereon the objective lens 2, the objective lens drive mechanism and an optical device for radiating a laser beam to a disc D through the objective lens 2. The base 1 is disposed between a guide shaft 3 and a screw shaft 4. The guide shaft 3 and the screw shaft 4 are arranged in an opposed parallel relation to each other on a mechanical chassis (not shown). The screw shaft 4 is rotated in both forward and reverse directions by means of a thread motor (not shown). Bearing portions 1a and 1b as portions to be guided are projected from both sides of the base 1. One bearing portion 1a is in abutment against a peripheral surface of the guide shaft 3, while the other bearing portions 1b are loosely fitted on the screw shaft 4. A plate spring 5 is screwed in a cantilevered fashion to the base 1 and a female screw member 6 called half nut is fixed to a free end side of the plate spring 5. The female screw member 6 is toothed over an approximately semicircle. This toothed portion remains in contact with the threaded portion of the screw shaft 4 due to the resilience of the plate spring 5.\nIn the optical pickup transfer mechanism thus roughly constructed, when the screw shaft 4 is rotated in either forward or reverse direction, the rotational force thereof is converted to a linear motion and transferred to the base 1 by means of the female screw member 6, so that the base 1 moves in the thrust direction (right and left direction in FIG. 7) of the guide shaft 3 and the screw shaft 4. As a result, the whole of the optical pickup 2 is moved radially along the disc D through the base 1, thus permitting information recording and reproducing operations for the disc.\nRecently, for promoting the reduction of cost and weight of the optical pickup moving mechanism, studies have been made about a technique wherein the base of the optical pickup is formed by molding a synthetic resin instead of using such a metallic material as die casting aluminum, and the guide shaft is formed by molding a synthetic resin as a substitute for a metallic shaft such as a stainless steel shaft. If, however, the base and the guide shaft are each formed by molding a synthetic resin, the bearing portion of the base, which is always kept in sliding contact with the guide shaft at the same portion thereof, becomes worn-out. This may result in tilting of an optical axis of the objective lens mounted on the base, which deteriorates the recording or reproducing function.\nIn the case of a vehicular disc unit, it is necessary to select a material superior in both rigidity and heat resistance. With a synthetic resin material, PPS (polyphenylene sulfide) with glass fibers incorporated therein is used. If both base and guide shaft repeat sliding, their constituent synthetic resins are apt to wear more rapidly. Such a problem is also true of the case where the screw shaft as the other movement guide member is formed by molding a synthetic resin.\nThe present invention has been accomplished in view of the above-mentioned circumstances and it is an object of the invention to provide an optical pickup moving mechanism wherein even if a base of an optical pickup and a movement guide member are each formed from a synthetic resin, it is possible to suppress the growth of wear caused by sliding contact between their constituent synthetic resins, and which can attain the reduction of cost without impairing the reliability.\nAccording to the present invention, for achieving the above-mentioned object, there is provided an optical pickup moving mechanism including an optical pickup and a movement guide member, the optical pickup having a base with an objective lens mounted thereon and with a to be guided portion being projected sideways thereof, the movement guide member having a guide surface which is extended in a predetermined direction and supporting the base slidably through the to be guided portion, the base being slid along the guide surface of the movement guide member, thereby allowing the optical pickup to move radially of a disc, wherein the to be guided portion and the movement guide member are each formed from a synthetic resin and a metallic slide member is disposed in the to be guided portion at a position opposed to the guide surface so that an outer peripheral surface of the slide member comes into sliding contact with the guide surface.\nAccording to the optical pickup moving mechanism of the above construction, since the base which includes the to be guided portion and the movement guide member which supports the base slidably are each formed by molding a synthetic resin, it is possible to attain the reduction of cost and weight, and since the metallic slide member disposed in the to be guided portion is brought into sliding contact with the guide surface of the movement guide member, it is possible to suppress the growth of wear caused by sliding contact between their constituent synthetic resins and hence possible to ensure a high reliability.\nIn the above optical pickup moving mechanism, by disposing a resilient member for urging the to be guided portion in a direction in which the slide member comes into pressure contact with the guide surface of the movement guide member, the slide member is allowed to slide on the guide surface always under a moderate pressing force without leaving the guide surface.\nIt is preferable that the slide member be a metal pin of a generally circular section. As such a metal pin of a generally circular section there may be used, for example, a commercially available parallel pin made of stainless steel which is inexpensive and high in dimensional accuracy. Thus, it is not necessary to newly provide a special mold for the metal pin; in other words, the metal pin itself does not become a cause of increase in cost. If the metal pin of a generally circular section is reciprocated in the optical pickup moving direction while being kept in sliding contact with the guide surface of the movement guide member, lubricating oil (grease) applied to the guide surface stays in wedge-like gaps formed before and behind the sliding contact position of the metal pin, thereby functioning as grease sumps and thus giving rise to an advantage that a smooth sliding motion of the metal pin is ensured.\nIn connection with the above optical pickup moving mechanism, if there is adopted a structure wherein a recess is formed in the to be guided portion at a position opposed to the guide surface and the metal pin is fitted in the recess, the metal pin mounting workability for the to be guided portion is improved. Further, if the recess is formed with a flat inner bottom surface and is narrower on its side close to the guide surface and if the metal pin is fitted in the recess while its outer peripheral surface is brought into abutment against the inner bottom surface of the recess, the metal pin can be fitted in the recess with the inner bottom surface of the recess as a reference. Here, the position of the bottom surface can be easily defined relative to the guide surface. This allows for improved relative positional accuracy between the metal pin and the guide surface and easier dimensional management in a design stage or in an assembling stage.\nIn connection with the above optical pickup transfer mechanism, if a through hole is formed in the to be guided portion so as to communicate with the recess and open to a side of the metal pin which side is not opposed to the guide surface, and if an adhesive is poured from this through hole into the recess to fix the metal pin within the recess, the metal pin fitted in the recess can be fixed firmly to the to be guided portion with the adhesive by using a slight force. Additionally, the assembling work efficiency can be improved because there is no fear that the adhesive may adhere to the guide surface side of the metal pin.\nPreferably, the metal pin is disposed so that its longitudinal direction is substantially orthogonal to the extending direction of the guide surface. According to this arrangement, the base can be moved while keeping the contact area between the metal pin and the guide surface to a minimum, so that the sliding resistance between the metal pin and the guide surface becomes extremely small and hence it is possible to prevent the occurrence of such an inconvenience as the metal pin is caught on the guide surface during movement of the optical pickup and tilting of the optical pickup results. For example, if the guide surface of the movement guide member is arcuate in section, then by disposing a metal pin in a direction approximately orthogonal to a generator of the circular arc it is made possible to maintain the metal pin and the guide surface in a state of point contact."}
{"text": "Powder coatings, which are dry, finely divided, free flowing, solid materials at room temperature, have gained considerable popularity in recent years over liquid coatings for a number of reasons. For one, powder coatings are user and environmentally friendly materials, since they are virtually free of harmful fugitive organic solvent carriers that are normally present in liquid coatings. Powder coatings, therefore, give off little, if any, volatile materials to the environment when cured. This eliminates the solvent emission problems associated with liquid coatings, such as air pollution and dangers to the health of workers employed in coating operations.\nPowder coatings are also clean and convenient to use. They are applied in a clean manner over the substrate, since they are in dry, solid form. The powders are easily swept up in the event of a spill and do not require special cleaning and spill containment supplies, as do liquid coatings. Working hygiene is, thus, improved. No messy liquids are used that adhere to worker's clothes and to the coating equipment, which leads to increased machine downtime and clean up costs.\nPowder coatings are essentially 100% recyclable. Over sprayed powders can be fully reclaimed and recombined with the powder feed. This provides very high coating efficiencies and also substantially reduces the amount of waste generated. Recycling of liquid coatings during application is not done, which leads to increased waste and hazardous waste disposal costs.\nIn the past, most powder coating was performed on metals which can withstand high temperatures at which many conventional coating powders fuse and cure. Recently, however, several coating powders have been developed for substrates, such as wood, which require coating powders which fuse (in the case of thermoplastic coating powders) or fuse and cure (in the case of curable coating powders) at relatively low temperatures. Examples of such coating powders are found, for example, in U.S. Pat. Nos. 5,824,373, 5,714,206, 5,721,052, and 5,731,043, the teachings of each of which are incorporated herein by reference. Low temperature coating prevents charring of the substrate and excessive outgassing of moisture.\nA frequent problem encountered when coating low-temperature substrates, such as wood, with coating powder is non-uniformity of coating in areas of the substrate which are difficult to coat, such as the edges and corners of kitchen cabinet doors. It has been found that preheating wood substrates, particularly in the 200.degree. F. to 275.degree. F. range, prior to electrostatic application of coating powders, provides more uniform coating of flat surfaces but can dry out sharp edges, making electrostatic coating difficult.\nHowever, preheating to near or above the boiling point of water tends to dry cellulosic products such as wood, fiberboard, particle board, paper, etc. Such materials tend to have a residual water content, wood typically having a water content of between about 3 and about 10 wt %. This residual moisture presents problems in coating cellulosic substrates with coating powder in that if the temperature is too high, significant outgassing causes defects, e.g., pinholes, in the coating. Similar problems have been noticed with fiber-containing plastic. This is one reason why cellulosic substrates must be coated with powders that fuse and cure at relatively low temperatures. On the other hand, the residual moisture in cellulosic materials is necessary for the material to hold sufficient electrical charge to be electrostatically coated with coating powder. Preheating of cellulosic substrates for the purpose of achieving uniform, continuous coatings may reduce the water content to where the charge-carrying capacity of the substrate is so reduced that electrostatic application of the coating powder is inefficient. Accordingly, it is a general object of the present invention to be able to preheat a substrate, such as a cellulosic substrate, for the purpose of achieving a uniform, continuous coating and at the same time maintaining sufficient moisture level of the substrate for electrical charge-carrying purposes.\nIn this regard, it was proposed to moisten the surface of lignocellulosic substrates prior to the pre-heating step so that the substrates would retain sufficient moisture and charge-carrying capacity at the point of electrostatic coating powder application. An example of this approach is found in above-referenced U.S. Pat. No. 5,824,373 which teaches maintaining substrates in a high humidity environment prior to pre-heating and optional humidity control through application of the coating powder. This approach, however, was discarded because it tended to warp the surface of the substrate."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates to a bladed rotor, and more particularly relates to a bladed rotor for a turbo-machine such as a gas turbine engine. The disclosure is particularly suited for use in gas turbine compressor rotors, although it is to be appreciated that the disclosure is not limited to compressor rotors and could find application in other types of bladed rotors for use in other types of turbo-machines.\n2. Description of the Related Art\nConventional axial compressor rotors for gas turbine engines typically comprise a number of discs which are bolted or welded together to form an integral rotatable drum. Each disc can be considered to represent a central hub around which a plurality of rotor blades of aerofoil configuration are mounted. Each rotor blade is normally attached to the hub using a mechanical connection known as a root fixing. One such type of arrangement involves axially fixing the rotor blades to the periphery of the hub and involves the provision of a series of slots which are machined into the peripheral region of the hub and which are generally elongate parallel to one another. The slots are typically arranged so that they extend in a lengthwise direction which makes an acute angle of between 10 and 30 degrees to the rotational axis of the hub. Each slot is configured to receive a dove-tail or fir-tree shaped root fixing of a respective rotor blade.\nA radially outwardly biased sprung retaining ring is normally used to secure the root portions of the rotor blades within their respective mounting slots. The retention ring locates within radially inwardly open grooves formed around the hub at positions located between the blade mounting slots, under its radially outward bias. Similar grooves are provided on the rotor blades and so the retaining ring also locates in the blade grooves to axially retain the root portions of the blades in the mounting slots.\nIt is important for integrity reasons that during operation of the rotor that the retaining ring does not apply radial load to the blades within the blade grooves. The retaining ring must at all times remain radially inwardly spaced from the radially outmost region of each blade groove by a clearance gap. It is therefore normal to configure the arrangement such that the retaining ring only bears against the radially outmost regions of the hub grooves.\nHowever, it has been found that during service the retaining rings of the above-described type of axial fixing arrangement can be susceptible to wear on their radially outmost surfaces, as also can the inner surfaces of the hub grooves within which the rings locate. Over time, this wear can reduce the size of the radial clearance gap between the retaining ring and the blade grooves which, as indicated above, cannot be allowed to occur due to integrity concerns."}
{"text": "In recent years, systems for playing back the contents of optical discs such as video CDs that record video data, audio data, and the like have been developed, and have prevailed for the purpose of playing back movie software titles, karaoke data, and the like.\nAmong such systems, a DVD (Digital Versatile Disc) standard that uses MPEG2 (Moving Picture Experts Group 2) international standards has been proposed.\nThis standard supports MPEG2 as a moving picture compression scheme, and AC-3 audio, MPEG audio, and the like as audio schemes. The standard is appended with sub-picture data for superimposed dialogs and menus obtained by compressing bitmap data, and control data (navigation data) for special playback control. Furthermore, this standard supports the UDF (Universal Disc Format) Bridge (a hybrid of UDF and ISO9660) to allow a computer to read data.\nAlso, optical discs such as a DVD-RAM and the like on which digital data can be written or rewritten have been developed. A function that allows one to easily edit the recorded contents is required of a digital video system using such DVD-RAM or the like.\nHowever, a home-use digital video system which allows end-users to easily edit the recorded contents has not become available yet."}
{"text": "1. Field of the Invention\nThe present invention relates to a signal detecting apparatus in an optical disc apparatus, in which a tracking error is detected by a push-pull method and a focusing error is detected by a spot-size method.\n2. Description of Related Art\nIn a known optical disc apparatus, a bundle of rays reflected by an optical disc is received by light receiving surfaces of error detecting light receiving elements. The signal data output from the divided areas of the light receiving elements is used to calculate the tracking error and focusing error.\nFor instance, Japanese Unexamined Patent Publication (kokai) No. 61-206944 (JPP '944) (U.S. Pat. No. 4,742,218) discloses a focus error detecting system in which light reflected from an optical disc is convered into a beam spot so that the focus state of an objective lens can be detected in accordance with the size of the beam spot. This known error detecting method will be referred to as a spot-size method hereinafter.\nIn the spot-size method, light receiving elements are located on opposite sides of and optically equidistant from a convergence point on which the light reflected from the optical disc is converged when the objective lens is in a focused state, so that the sizes of the respective beam spots formed on the light receiving elements can be compared with each other to generate a focus error signal.\nHowever, in an optical system as disclosed in FIG. 1 of JPP '944 mentioned above, it is impossible to detect tracking error with the same optical elements that are used for detecting the focusing error. In the optical system shown in FIG. 4 of JPP '944, both tracking error and focusing error can be detected with the same optical elements. However, the pattern into which the light receiving elements are divided is complex.\nFurthermore, if the optical axis of the bundle of rays, incident upon the objective lens, is inclined with respect to the optical disc, or if the objective lens is displaced in the radial direction of the optical disc to correct tracking error, the optical path of the reflected light is deviated from a reference position. This results in the displacement of the beam spots in a direction corresponding to the radial direction of the optical disc. Consequently, in a tracking error detecting system that uses the push-pull method, even if there is no change in light intensity distribution, which is caused when the beam spot moves across the optical disc, the light receiving areas will be unbalanced, resulting in track offset signals being carried on the detected signals. Therefore, the deviation of the beam spots from the track of the optical disc and the track error signals no longer maintain a predetermined relationship, and accordingly, the position of the beam spot can not be precisely controlled by a track servo control in accordance with the detected signals.\nThe \"track offset signal\" is one of the track error signal components detected by the light receiving elements. Track offset is caused by the displacement of the beam spot on the light receiving element due to the deviation of the reflected light.\nIn the known focusing error detection system using the spot-size method, the resultant signal of the outputs of the two light receiving elements is set to be zero when the objective lens is in the focal position. Namely, the respective signals of the light receiving elements are not individually taken into account.\nConsequently, it is necessary to adjust the position of the light receiving elements while observing the balance of the quantities of light to be received by the light receiving elements, thus resulting in a complex and troublesome adjustment operation.\nFurthermore, if the light receiving elements are used to reproduce the recorded magnetic optical signals of the magnetic optical disc as disclosed in the above-mentioned JPP '944, a slight change in the balance of the quantity of light to be received by the light receiving elements due to the rotation of the polarizing surface by the Kerr effect occurs, and accordingly, the change may result in an interference contained in the focus error signal, thus resulting in an imprecise focus servo control.\nIn the prior art, as disclosed in JPP '944, in which each light receiving element is split into three light receiving areas, the tracking error, using the push-pull method, and the focusing error cannot be detected by the same light receiving element.\nFurthermore, JPP '944 also discloses light receiving elements, each being split into three light receiving sections in the form of elongated bands, wherein the center light receiving section is split into three mosaic areas, so that the tracking error can be detected by the same light receiving element as that for detecting the focusing error. The split pattern of the light receiving element is, however, complex, especially at the center portion thereof. Accordingly, it is necessary to form a relatively large beam spot on the light receiving element, thus resulting in a decreased freedom of optical design and a large apparatus."}
{"text": "1. Field of the Invention\nThis invention relates to a method for forming a pattern of a film made of a desired material on a side wall of a stepped base, which pattern formation method is appropriately used, for example, for forming a desired film pattern of good precision on a plane which extends orthogonally to the main plane of the base in order to achieve high integration of an LSI.\n2. Description of the Related Art\nIn order to improve the degree of integration of an LSI over a limited substrate area, the reduction of areas over which individual semiconductor elements constituting the LSI are defined on the main plane of the substrate has become more and more necessary. To this end, a technique has been employed wherein a groove is made in the substrate so that the side walls (portions corresponding to wall surfaces) are utilized as regions for forming semiconductor elements in order to provide an increased surface area over which the elements can be formed. However, the side walls defining such grooves have only been used, for example, to accommodate capacitor elements and have never been used to form a wiring pattern or a gate electrode of a transistor. Nevertheless, the necessity for forming a wiring pattern or elements on the side wall will be more and more in demand, and so it is believed to be necessary to establish a method wherein a wiring pattern or a gate electrode can be readily formed.\nFIG. 5 illustrates a thin film pattern on a side wall, and is a perspective view of an ideal form of a thin film pattern. In this case, a semiconductor substrate 13 has a groove 11 with a depth l, and the thin film pattern 17 extends on a side wall 15 defining the groove 11 along the depth l of the groove 11.\nHowever, it is very difficult to form such an ideal thin film pattern 17 as shown in FIG. 5.\nFIGS. 6(A)-(C) illustrate main steps used to form a film pattern on a side wall according to a pattern forming method using the known single layer resist process. More particularly, FIGS. 6(A) and (C) are, respectively, partial perspective views taken in the direction P shown in FIG. 5, and FIG. 6(B) is a plan view of the sample shown in FIG. 6(A).\nWhen the single resist process is used, a film 21 of material used to form a thin film pattern is formed on the surface of the semiconductor substrate 13 including over the inner surfaces thereof defining the groove 11. Next, the film material 21 is coated over the entire surface thereof, for example, with a positive resist (not shown). Subsequently, a light-shielding mask is provided at such a position that it crosses a boundary 11a of the step established by the groove 11, the resist is exposed to light irradiated from above the light-shielding mask, and the resist is developed to form a resist pattern 25 on a region traversing the boundary 11a of the step (FIG. 6(A)). However, the thus formed resist pattern 25 has a shape which is completely different from the shape of the light-shielding mask 23 as is particularly shown in FIG. 6(B). The reason for this is that the resist which has been applied over the surface of the semiconductor substrate 13 is thicker at the bottom portion of the groove and particularly at the corner portions of the groove 11, so that the exposure light does not reach the lower side of the thicker portion of the resist. Accordingly, a thin film pattern 21a which is obtained by an anisotropic etching technique, such as RIE (Reactive Ion Etching), has a shape which is far different from the shape of the light-shielding mask 23 as is shown in FIG. 6(C). More particularly, a useless region as shown in FIG. 6(C) as hatched is left at the lower portion of the step.\nAs a measure for overcoming the drawback involved in the known single layer resist process, there was known a so-called double layer resist process disclosed, for example, in \"Process Techniques For Next Generation Super LSI,\" Applications, Apr. 4, 1988, by Realize Co., Ltd., (pp. 297-298). In this process, after the step of the substrate is filled in with a first resist layer, a second resist layer is formed on the first resist layer and is exposed to light and developed to obtain a mask. The first resist layer is subjected to patterning through the mask, thereby precisely controlling the patterned shape at the stepped portion.\nFIGS. 7(A)-(F) illustrate the principle of the double layer resist process wherein there are shown steps of forming a film pattern on side walls 31a, 33a of two recesses 31, 33 of a semiconductor substrate 35, respectively. These figures are sectional views, respectively, of the semiconductor substrate 35 taken along the direction orthogonal to the side walls 31a, 33a.\nInitially, a film 37 of material is formed on the entire surface of the semiconductor substrate 35 having the recesses 31, 33. A first thick resist layer 39 is then provided to define a flat upper surface (FIG. 7(A)).\nThen, a second resist layer 41 is formed on the first resist layer 39 (FIG. 7(B)), after which the second resist layer 41 is subjected to exposure light in a desired pattern and developed to form resist patterns 41a, 41b (FIG. 7(C)).\nSubsequently, the first resist layer 39 is subjected to patterning by, for example, dry etching to form the double layer resist patterns 43a, 43b (FIG. 7(D)).\nThe double layer resist patterns 43a, 43b are used as a mask for patterning the film material 37 so that portions 37a, 37b of the film material 37 are left over the respective side walls 31a, 33a defining the recesses 31, 33 (FIG. 7(E)).\nThen, the double layer resist patterns 43a, 43b are removed to expose the desired film patterns 37a, 37b (FIG. 7(F)).\nAccording to the double layer resist process, the flattening layer (first resist layer) and the layer exposed to light and subjected to development (second resist layer) are provided separately, so that a pattern shape corresponding to the shape of the exposure mask can be obtained even at the stepped portion.\nHowever, the double layer resist process is complicated. As will be apparent from the process illustrated with reference to FIGS. 7(A)-(F), the double layer resist process has additional steps, including the step of applying the first resist layer and the step of patterning the first resist layer, compared with the single layer resist process. Moreover, although omitted in the above description, the double layer resist process inevitably requires a baking step after the application of the first resist layer and prior to the application of the second resist layer, wherein heat is applied to dry the first resist layer.\nIn addition, when the first resist layer is anisotropically etched, the second resist layer should have a satisfactory etching selection ratio to the first resist layer, thus presenting a problem in that the selection of materials is difficult.\nThe complicated process and the reduction in degree of freedom with respect to the selection of resist materials are not favorable, for example, in view of the fact that they contribute to a direct rise in the production costs of the semiconductor device.\nBoth the known single layer resist process and the known double layer resist process are disadvantageous in that the film pattern cannot be left only on the side wall or walls of the step (e.g. the ideal film pattern 17 should be left only on the side wall as shown in FIG. 5), i.e. the film pattern also remains on the upper and lower portions of the step. This is because in the prior processes, after the film material has been formed on the entire surface of the substrate including the step, a resist pattern is formed on the film material so as to cover the side wall. More particularly, in order to reliably cover the side wall with the resist pattern, a shift in the alignment between the mask used to form the resist pattern (the light-shielding mask as in FIG. 6(A)) and the side wall has to be taken into account. Accordingly, the light-shielding mask 23 must inevitably extend over the upper and lower portions of the step while covering the step. Since the material film 21 is formed on the entire surface of the substrate, the film pattern 21a will be left on the upper and lower portions of the step as corresponding to the additional regions required for the mask alignment. This phenomenon takes place in the double layer resist process illustrated with reference to FIG. 7."}
{"text": "In recent years, numerous improvements have been made in eyewear design and production, which have produced increasingly lightweight, comfortable, and attractive products. However, these product improvements have exposed a number of problems with traditional eyewear designs. First, many eyewear products are very delicate, have multiple small components, and require sophisticated tools or custom parts for repair or assembly. Second, the delicate components inevitably break after extended use, and the repairs are time consuming, expensive, or inconvenient for the user. Third, for a given set of eyeglasses frames, it is often difficult, expensive, or impossible to customize the eyeglasses for a given user. In addition, repairs to damaged eyeglasses can often not be made by the end user, necessitating a visit to an optician. The present invention can provide a solution to at least one of these problems.\nA major problem with many eyeglasses designs is the need for multiple fasteners such as screws, pins, or small bolts. These fasteners may be located at hinge points between the eyeglasses temple arm and temple arm hinge block or at various positions on the eyeglasses rim. In some cases, multiple fastener types or sizes are used on the same pair of eyeglasses. Furthermore, these fasteners are often not easily replaced through a vender's stock and often require custom orders to make repairs.\nFasteners may be used both on eyeglasses rims and eyeglasses hinges. When used on eyeglasses rims, the fasteners may help secure a lens in place, connect the rim to another portion of the eyeglass, or hold multi-component parts together. In any case, the fastener, whether it be a screw, pin, or bolt design, is subject to back-out after extended wear. When this occurs, the fastener will often be lost, and the product will be rendered useless until a repair is made.\nFasteners are also used for eyeglasses hinges, for which numerous designs exist. Some designs require threaded fasteners such as screws or small bolts. Others require threadless fasteners such as pins. Some designs also employ an adhesive, washers, or friction-fit materials. Regardless, with all of these hinge designs, the screw, bolt, or pin risks backing out of the socket in the hinge or other frame section, rendering the eyeglasses unwearable and potentially requiring the purchase of new parts, the use of special tools, or a consultation with an optician to make repairs.\nAnother problem with eyeglasses hinges is that they are sometimes subject to relatively severe stress due to accidental or intentional misuse. Traditional eyeglasses hinges will often break or become distorted under sufficient stress. Broken eyeglasses cannot be worn, and distorted eyeglasses may fit improperly. With most current designs, repairs may require significant training, the purchase of replacement parts, or the use of unsightly materials such as tape or glue.\nYet another problem with current eyeglasses designs pertains to the temple arm ear piece. The temple arm ear piece comes in one length for a given set of frames. Although consumers often need shorter or longer temple arm lengths, this part is rarely stocked and must be special ordered by a vendor. This process is time-consuming if the appropriate piece is available at all. Most temple arms on frames are not universal, are typically left and right-sided, and cannot be switched with other frame styles. The earpiece sock or paddle is not designed to be removed once applied by the manufacturer. In the event that the paddle is lost or breaks, replacing it is very difficult. Paddle styles are often custom items designed for only one temple arm style, color, and size; and manufacturers do not often use paddle part numbers for reordering. Most provide complete frames only for replacement."}
{"text": "Generally described, computing devices utilize a communication network, or a series of communication networks, to exchange data. Companies and organizations operate computer networks that interconnect a number of computing devices to support operations or provide services to third parties. The computing systems can be located in a single geographic location or located in multiple, distinct geographic locations (e.g., interconnected via private or public communication networks). Specifically, data centers or data processing centers, herein generally referred to as a “data center,” may include a number of interconnected computing systems to provide computing resources to users of the data center. The data centers may be private data centers operated on behalf of an organization or public data centers operated on behalf, or for the benefit of, the general public.\nTo facilitate increased utilization of data center resources, virtualization technologies may allow a single physical computing device to host one or more instances of virtual machines that appear and operate as independent computing devices to users of a data center. With virtualization, a single physical computing device can create, maintain, delete, or otherwise manage virtual machines in a dynamic matter. In the simplest embodiment, users can request single computing device computer resources from a data center. In more complex embodiments, users, such as system administrators, can request the configuration of virtual machine instances corresponding to a desired set of networked computing devices. In such embodiments, the data center can implement varying number of virtual machine instances to implement the functionality and configuration of the requested physical computing device network, generally referred to as a virtual machine network.\nFor virtual machine network embodiments, users often want to utilize various services, components (such as network-based appliances), or other functionality in accordance with at least aspects of the implementation of a hosted virtual machine network. In one aspect, users are required to configure various information about a hosted network, such as address space information, domain name service (DNS) zones information, resilient packet transport (RPT) information, and the like, in order for the desired functionality to be implemented in the hosted virtual machine network. Additionally, in another aspect, users are also required to delegate access to, or otherwise grant access, to at least a portion of the hosted virtual machine network to the entity providing desired functionality. As such, the virtual machine network service provider would prefer for users to be aware of the type permissions or authorizations that are delegated in conjunction with the utilization of requested functionality. Current approaches to the management of configuration information and delegated permission information are ad hoc in nature."}
{"text": "In the past, music games have been enjoyed in which the progress of a game is controlled by causing an indication sign having a shape of a musical instrument such as a drum to move along a predetermined path on a display screen so as to be matched with a piece of music, determining a propriety of an operation input when the indication sign reaches a predetermined criterion position, and adding a score (for example, see JP-A-2005-87323). In most of such music games, when an operation input to the indication sign could not be performed, a part of a piece of music to be output is not output merely instead of changing the piece of music."}
{"text": "This invention relates to an optical connector and, more particularly, to an optical connector incorporating a polarization-independent optical isolator to be provided between optical fibers in fiber optics communication systems and the like.\nWith the recent advances in optical communications that use semiconductor lasers as signal light sources, it has become possible to transmit signals at high speed and density exceeding several gigahertz. Among the various optical components used in such high-speed and density signal transmission is an optical isolator which prevents the reentrance of reflected light into semiconductor lasers.\nOptical isolators are of two types, a polarization-dependent isolator which transmits only the light travelling in a specified direction of polarization and a polarization-independent isolator which transmits light in any direction of polarization. The second type of optical isolators are typically used in light amplifiers at repeaters in signal transmission systems and will be in great demand in the future.\nFIG. 26 shows the construction of a typical example of the conventional polarization-independent optical isolator. Generally indicated by 410, the isolator comprises one Faraday rotator and three birefringent crystal plates.\nIn FIG. 26, the first to the third birefringent crystal plates are identified by 411, 412 and 413, respectively, and the Faraday rotator 414 is provided between the plates 411 and 412. A magnetic field parallel to the Z direction is applied to the Faraday rotator 414. The birefringent crystal plates 411, 412 and 413 are parallel plates prepared by polishing slices of a uniaxial crystal that have been cut in such a way that their C axis is at an angle with the surface. A ray of light incident normal to each of these parallel plates is separated into two components that are polarized in orthogonal crossed directions. The birefringent crystal plates 411, 412 and 413 have different thicknesses in the direction of light transmission and their ratio is 1:1/.sqroot.2:1/.sqroot.2. The plate 413 is such that its C axis coincides with the C axis of the plate 412 if the latter is rotated through 90.degree. about the Z axis. The Faraday rotator 414 is typically formed of a bismuth-substituted garnet and can rotate the direction of light polarization non-reciprocally through an angle of 45.degree.. Shown by 415 is a coupling lens for coupling the light to an optical fiber 416 or 417.\nFor the purpose of the following discussion, the direction of light travel is assumed to be \"forward\" if it is launched from the birefringent crystal plate 411 and \"backward\" if it is launched from the plate 413. Thus, the forward incident ray of light is indicated by 410f and the backward incident ray of light is indicated by 410b. When the incident light is separated into two components, those in the forward direction are indicated by f1 and f2 whereas those in the backward direction are indicated by b1 and b2. The direction of light travel is represented by the arrow.\nFIG. 27 shows how light travels through the optical isolator when it is seen from the birefringent crystal plate 411. A part (1) of FIG. 27 refers to the case of forward light propagation and a part (2) of FIG. 27 refers to the case of backward light propagation; A-E correspond to the respective positions A-E in FIG. 26; the dots represent the positions of respective light components and the arrows represent the directions of planes of polarization. The plane of polarization is assumed to rotate in \"+\" direction if it rotates clockwise.\nThe operating principle of the optical isolator will now be described with reference to FIGS. 26 and 27. If the C axis of the birefringent crystal plate 411 is directed upward (along the Y axis), forward signal light 410f launched from the coupling lens 415 to be incident on the plate 411 is separated into two components f1 and f2 in orthogonal crossed directions of polarization (see at B in FIG. 27(1)). With their relative positions remaining the same, the components f1 and f2 have the respective planes of polarization rotated through +45 degrees by the Faraday rotator 414 and then enter the birefringent crystal plate 412 (see at C in FIG. 27(1)). The plate 412 is such that its C axis coincides with the C axis of the birefringent crystal plate 411 if the latter is rotated through -45 degrees, so the component f1 is refracted as an extraordinary component whereas the component f2 which is an ordinary component is not diffracted but simply transmitted through the plate 412 (see at D in FIG. 27(1)). The birefringent crystal plate 413 is such that its C axis coincides with the C axis of the plate 412 if the latter is rotated through +90 degrees, so the component f2 is refracted as an extraordinary component whereas the component f1 which is an ordinary component is simply transmitted through the plate 413 (see at E in FIG. 27(1)). Thus, the two components of polarization are recombined at point E and coupled to the optical fiber 416 by means of the coupling lens 415.\nThe backward light 410b, as far as it travels to point C, behaves in essentially the same way as the forward light 410f, except that due-to the non-reciprocity of the Faraday rotator 414, the incident light components b1 and b2 have their planes of polarization rotated through +45 degrees as seen in the forward direction before they are incident on the birefringent crystal plate 411 (see at B in FIG. 27(2)). As a result, the component b1 is refracted as an extraordinary component whereas the component b2 which is an ordinary component is simply transmitted through the plate 411 (see at A in FIG. 27(2)). Thus, the components b1 and b2 emerge from the birefringent crystal plate 411 in different positions than when the forward light was launched into the same plate 411 and, hence, they will not couple with the optical fiber 417, thereby insuring that the reflected light will be isolated from the semiconductor laser.\nFIG. 28 shows the exterior appearance of a conventional polarization-independent optical isolator. The optical isolator generally indicated by 420 comprises an isolator portion 418 and a connector portion 419 at both ends. The isolator portion 418 has the components that are shown in FIG. 26 and which are adjusted and fixed within a case. The connectors 419 are connected to optical fibers in other transmission systems. The size of the optical isolator portion 18 may be about 7 mm in diameter and 45 mm long.\nThe conventional polarization-independent optical isolator comprising a plurality of birefringent polarizing plates and a single Faraday rotator has suffered from the following disadvantages.\n(1) It contains many parts that need precise optical adjustments, so the number of steps involved in assembly is so great as to make it a cumbersome and time-consuming operation.\n(2) When the optical isolator portion is to be coupled to optical fibers, the great number of its components increases the length of the space through which light propagates between fibers. In addition, the rays of light that entered in the forward direction will emerge at positions deviating from the axes of the incident rays; therefore, the positions that serve as guides for the coupling lenses and optical fibers at opposite ends cannot be uniquely determined and, hence, considerable labor is needed to achieve optical axial alignment.\n(3) Coupling to other transmission systems is only accomplished by means of the connectors at opposite ends, so a large installation space is required to incorporate the optical isolator into a measuring instrument or communication equipment.\n(4) The individual optical devices are provided normal to the optical fibers, so the reflected light from these optical devices will return to the optical fibers such as to deteriorate the system's reflection attenuation characteristics.\n(5) If the optical isolator is to be provided on the exit side of an optical fiber amplifier, a separate wavelength filter operating over a narrow band of frequencies is necessary but then the construction of the amplifier becomes complicated."}
{"text": "Sand consolidation is a well known term applying to procedures routinely practiced in the commercial production of petroleum, whereby wells are treated in order to reduce a problem generally referred to as unconsolidated sand production. When wells are completed in petroleum-containing formations which also contain unconsolidated granular mineral material such as sand or gravel, production of fluids from the formation causes the flow of the particulate matter into the wellbore, which often leads to any of several difficult and expensive problems. Sometimes a well is said to \"sand up\", meaning the lower portion of the production well becomes filled with sand, after which further production of fluid from the formation becomes difficult or impossible. In other instances, sand production along with the fluid results in passage of granular mineral material into the pump and associated hardware of the producing well, which causes accelerated wear of the mechanical components of the producing oil well. Sustained production of sand sometimes forms a cavity in the formation which collapses and destroys the well. All of these problems are known to exist and many methods have been disclosed in the prior art and applied in oil fields in order to reduce or eliminate production of unconsolidated sand from a petroleum formation during the course of oil production.\nThe above-described problem and potential solutions to the problem have been the subject of extensive research by the petroleum industry in the hope of developing techniques which minimize or eliminate the movement of sand particles into the producing well and associated equipment during the course of producing fluids from the formation. One such general approach suggested in the prior art involves treating the porous, unconsolidated sand mass around the wellbore in order to cement the loose sand grains together, thereby forming a permeable consolidated sand mass which will allow production of fluids but which will restrain the movement of sand particles into the wellbore. The objective of such procedures is to create a permeable barrier or sieve adjacent to the perforations or other openings in the well casing which establish communication between the production formation and the production tubing, which restrains the flow of loose particulate mineral matter such as sand. Another approach involves removing a portion of the formation around the well and packing specially prepared resin-coated granular material into the formation around the wellbore which is subsequently caused to become cemented together.\nIt is a primary objective of any operable sand consolidation method to form a barrier around the wellbore which restrains the movement of sand particles into the well while offering little or no restriction to the flow of fluids, particularly oil, from the formation into the wellbore where it can be pumped to the surface of the earth.\nAnother very important quality of a satisfactory sand consolidation method is durability of the permeable barrier formed around the wellbore. Once the barrier is formed and the well is placed on production, there will be a substantial continuing flow of fluids through the flow channels within the permeable barrier, and it is important that the barrier last for a significant period of time, e.g. several months and preferably years, without excessive abrasive wear or other deterioration of the consolidation matrix which would allow the particulate matter to once again flow into the wellbore. This is a particularly difficult objective to accomplish in the instance of sand consolidation procedures applied to wells completed in formations subjected to steam flooding or other thermal recovery methods. The production of fluids in steam flooding operations involve higher temperatures and higher pH fluids than are normally encountered in ordinary primary production, and this greatly aggravates the stability problem of sand consolidation procedures.\nIt is also important that the material injected into the formation should be essentially unreactive during the period it is inside the wellbore, i.e. while it is being pumped down the well and positioned where it is desired adjacent to the perforations of the production casing. It is this desire to delay the consolidation reaction that has led to multi-step procedures in which first a catalyst is injected into the formation, after which the polymerizable resin-containing fluid is injected separately. While this reduces the propensity for the fluid to polymerize in the injection string, it does give rise to several problems which constitute inherent weaknesses in many prior art methods for accomplishing sand consolidation. First, each separate injection step increases the time and cost of the well treatment by which sand consolidation is accomplished. Second, when the only catalyst employed is injected into the formation in advance of the polymerizable fluid, uniform mixing of catalyst with all of the subsequently-injected polymerizable fluid is not achieved to the degree necessary to ensure optimum polymerization of the resin, and thus often fails to achieve maximum, uniform strength and durability of the consolidated mass.\nMany materials have been utilized for consolidating sand in the formation adjacent to production of wellbores. One of the more successful agents utilized for this purpose are fluids comprising monomers or oligomers of furfuryl alcohol which can be polymerized in situ to form a solid matrix which binds the sand grains together, while at the same time offering superior resistance to high temperatures and to caustic substances which may be encountered in steam flood operations. One of the problems in utilizing furfuryl alcohol oligomers to polymerize in the formation for the purpose of consolidating sand grains is failing to achieve uniform catalysis of the polymerization. Many catalysts that are effective for polymerizing furfuryl alcohol resins cannot be admixed with the furfuryl alcohol to permit a single fluid containing both the resin and the catalyst to be injected into the formation, because the time of polymerization is so short or unpredictable that there is excessive danger that the resin will polymerize in the injection wellbore.\nIn U.S. Pat. No. 4,427,069 there is disclosed a procedure for consolidating sand in a formation adjacent to a wellbore using an oligomer of furfuryl alcohol, in which the catalyst used is a water soluble acidic salt, preferably zirconyl chloride, which is injected in an aqueous solution into the formation prior to injection of the resin-containing fluid. The salt absorbs on the sand grains, and sufficient acidic salt remains adsorbed on the sand grain during the subsequent resin fluid injection stage that adequate polymerization occurs. Although this has been very effective in difficult situations where sand consolidation procedures are utilized, particularly in connection with thermal flooding such as steam injection procedures, the procedure nevertheless requires a multi-fluid injection procedure which requires more time and is more expensive than is desired. Usually a preliminary sand cleaning step is required before injecting the aqueous-catalyst solution in order to remove the naturally-occurring oil film from the sand grains to ensure good catalyst adsorption on the sand. Also, although catalyst mixes with the subsequently injected polymer to a limited degree, usually sufficient to cause polymerization, it is believed that superior performance would result if the catalyst resin mixture can be made more homogenous prior to polymerization, in order to achieve a dense, strong, durable consolidation mass.\nIn U.S. Pat. No. 4,669,543 which issued June 2, 1987, there is described a method for consolidating sand using an acid curable resin and utilizing as a catalyst, the reaction product of an acid and an alkyl metal or ammonium molybdate. In that instance, the catalyst is incorporated in an aqueous carrier fluid which comprises the continuous phase of an emulsion in which the polymerizable resin is the dispersed or discontinuous phase. Thus this process requires that the emulsion be resolved or broken after it is located in the portion of the formation where the permeable consolidating mass is desired, which is difficult to achieve the degree of completion and accuracy of timing necessary to accomplish the desired strong durable consolidating matrix necessary for a long lasting sand consolidation process.\nIn our U.S. Pat. No. 4,482,072 for \"SAND CONSOLIDATION METHODS\" we disclosed a particularly effective method for consolidating sand utilizing a mixture of a polymerizable resin such as an oligomer of furfuryl alcohol and a diluent such as butyl acetate and an oil soluble, slightly water soluble acid catalyst such as orthonitrobenzoic acid is injected followed by injection of salt water to reestablish permeability.\nIn U.S. Pat. No. 4,903,770 for \"SAND CONSOLIDATION\" we disclosed a preferred process which is more easily removed after a period of use and which is quite inexpensive. The process employs a fluid comprising a polymerizable monomer such as furfuryl alcohol and as a diluent, a polar organic solvent such as methanol and a strong, non-volatile acid catalyst such as sulfuric acid, mixed with steam to form a multiphase or aerosol treating fluid, and injected into the formation to be consolidated. An ester such as ethyl or butyl acetate is incorporated in the fluid when the steam quality is less than 80 percent.\nBoth of the above methods have produced excellent results in many field applications including several viscous oil-containing formations being stimulated by steam flooding. In some applications, however, results are adversely affected by a slight shrinkage of the porous, polymerized mass which can cause cracks in the polymerized mass which causes leaks. There is a need for an inexpensive modification to the sand control treatment using monomers or oligomers of furfuryl alcohol to reduce or eliminate shrinkage of the polymerized polymer-sand mass."}
{"text": "1. Field of the Invention\nThis invention relates in general to a parallel cipher text scrambler, and more particularly to cipher text scrambler and method for scrambling multiple, parallel data streams so that the data streams are uncorrelated and have a fixed delay therebetween.\n2. Description of Related Art\nCommunication information encryption is performed in digital data communications to provide security. Furthermore, communication applications require the sequence of data bits transmitted over a communication channel to be statistically random. In order to achieve the required degree of randomness, the data can be scrambled using a Maximal Length Pseudo Random Sequence. Maximal Length Pseudo Random Sequences are known to have the lowest possible auto-correlation, and are therefore the optimal choice for scrambling.\nTwo techniques have been generally used in the prior art to perform the scrambling operation: block scrambling and stream scrambling. Both techniques take advantage of the fact that when a first sequence of bits is exclusively OR'ed with a second sequence of bits and is then again exclusively OR'ed with the second sequence of bits identically aligned, the output is the first sequence just as it was before any exclusive OR operations were performed.\nBlock scrambling uses a framing pattern or other known means to provide the bits into some definable blocks of information. These bits are then exclusively OR'ed with a fixed pattern of bits synchronized to the boundaries of the block. Since a pattern of scrambling bits is fixed with respect to the block, the same pattern can be used at the receiver end to unscramble the bits. Any bit error occurring in the transmission channel between the transmitter and receiver will cause an error in that particular bit, but will not cause other bits to be in error, provided only that the receivers remain synchronized with the transmitters to the block boundaries. However, the longest framing pattern is not usually equal to the 2.sup.n-1 length of the pseudo random pattern desired, so that only part of the pattern is used repetitively, introducing undesirable correlations in the data.\nStream scrambling of the known prior art generally operates on a continuous stream of bits. In typical implementation, the bits at the transmitting end of a communication channel to be scrambled are passed through one input of a two input exclusive OR gate. The output of the gate is the output of the scrambler and also the input to an N-stage shift register. This shift register is tapped at the Nth stage and one or more other stages, and the outputs of these taps are exclusively OR'ed together. The result of this exclusive OR operation is applied to the other input of the exclusive OR gate that has the data to be scrambled, as the first input. The tap positions are chosen so that a Galois polynomial represented by the tap weight is irreducible, and if the input data were all zeros and the shift register started out at any state other than all zeros, a Maximal Length Pseudo Random Sequence would be produced.\nAs mentioned, pulse patterns can have an energy component which is particularly high at certain discrete frequencies. In order to avoid these pulse patterns, the digital signal to be transmitted must be scrambled at the transmitting side with a pseudo random sequence. The descrambling occurs at the receiving side with the pseudo random sequence which was employed at the transmitting side. The synchronization of the pseudo random generators employed at the transmitting receiving sides which is thereby necessary, can be avoided by employing freewheeling and, therefore, self synchronizing scrambler and descrambler arrangements. In order to achieve high bit rates, modulating signals having a high clock frequency are often scrambled in a plurality of parallel channels having a lower bit repetition frequency. However, the bit streams must be synchronized in order to reconstruct the original signal.\nThus it can be seen that there is a need for a cipher text scrambler which enables multiple data streams to be scrambled in parallel in such a way that the data streams are uncorrelated with respect to each other.\nIt can also be seen that there is a need for a cipher text scrambler that provides a large fixed delay relationship between the scrambler bits over the multiple streams.\nIt can also be seen that there is a need to provide a cipher text scrambler that can correctly identify the ordering of the multiple data streams with the respect to each other, correct for any polarity or wore swap misconnections and align the received multiple data streams despite differential delays between the data streams."}
{"text": "A user equipment (UE) may connect to a network via an access point, such as a base station, an eNodeB, a wireless local area network (WLAN) access point (e.g., a Wi-Fi access point), or the like. The UE may experience degraded network performance and/or may disconnect from the network when network congestion associated with a connection to the access point satisfies a network congestion threshold. The network congestion threshold may be satisfied when a threshold quantity of UEs attempt to utilize the same access point for network connectivity."}
{"text": "Autoimmune diseases, such as systemic lupus erythematosus (SLE), myasthenia gravis (MG) and idiopathic thrombocytopenic purpura (ITP), among others, remain clinically important diseases in humans. As the name implies, autoimmune diseases wreak their havoc through the body\"\"s own immune system. While the pathological mechanisms differ between individual types of autoimmune diseases, one general mechanism involves the binding of certain antibodies (referred to herein as self-reactive antibodies or autoantibodies) present in the sera of patients to self nuclear or cellular antigens.\nSLE has an incidence of about 1 in 700 women between the ages of 20 and 60. SLE can affect any organ system and can cause severe tissue damage. Numerous autoantibodies of differing specificity are present in SLE. SLE patients often produce autoantibodies having anti-DNA, anti-Ro, and anti-platelet specificity and which are capable of initiating clinical features of the disease, such as glomerulonephritis, arthritis, serositis, complete heart block in newborns, and hematologic abnormalities. These autoantibodies are also possibly related to central nervous system disturbances. Kidney damage, measured by the amount of proteinuria in the urine, is one of the most acute areas of damage associated with pathogenicity in SLE, and accounts for at least 50% of the mortality and morbidity of the disease. The presence of antibodies immunoreactive with double-stranded native DNA is used as a diagnostic marker for SLE.\nAntibodies are composed of heavy and light polypeptide chains which are joined by disulfide bridges. Antibodies are divided into different classes according to their heavy chain structure; antibodies belonging to the same class are referred to as isotypes of each other. In addition, antibodies of a given isotype can be divided into subtypes. Antigenic determinants on antibodies that differ among animals that have inherited different alleles are referred to as allotopes; antibodies that share an allotope are referred to as members of the same allotype. Another type of antigenic determinant present on antibody molecules are those found primarily in the hypervariable region of the antigen binding site of the antibody. These determinants are referred to as idiotopes; antibodies that share an idiotope are referred to as members of the same idiotype. Idiotypic determinants are controlled by both genetic and antigenic influences. Antibodies having common or shared idiotypes generally exhibit the same antigenic specificity. However, antibodies from genetically different individuals which share a common antigenic specificity may exhibit idiotypic heterogeneity but, in some instances, show a major cross-reactive antigenic determinant. Thus, antibodies which bind the same antigen may have distinct idiotypic determinants, but also may share cross-reacting properties.\nCurrently, there are no really curative treatments for patients that have been diagnosed with SLE. From a practical standpoint, physicians generally employ a number of powerful immunosuppressive drugs such as high-dose corticosteroids, azathioprine or cyclophosphamidexe2x80x94many of which have potentially harmful side effects to the patients being treated. In addition, these immunosuppressive drugs interfere with the person\"\"s ability to produce all antibodies, not just the self-reactive anti-DNA antibodies. Immunosuppressants also weaken the body\"\"s defense against other potential pathogens thereby making the patient extremely susceptible to infection and other potentially fatal diseases, such as cancer. In some of these instances, the side effects of current treatment modalities can be fatal.\nOne method of treatment for SLE, described in Diamond et al. (U.S. Pat No. 4,690,905), consists of generating monoclonal antibodies against anti-DNA antibodies (the monoclonal antibodies being referred to therein as anti-idiotypic antibodies) and then using these anti-idiotypic antibodies to remove the pathogenic anti-DNA antibodies from the patient\"\"s system. However, there are several drawbacks to this approach. For example, the removal of large quantities of blood for treatment can be a dangerous, complicated process. Essentially, blood is removed from a patient, treated to remove the anti-DNA antibodies, and then the treated blood returned to the patient. Such a removal technique would be similar to that used for hemodialysis, i.e., via an arterial passage. This type of treatment would be inconvenient (a qualified professional would be required to conduct treatment regularly), expensive, painful, and in some instances might subject the patient to a risk of infection and/or hemorrhaging, as well as depletion of effective blood volume inducing circulatory collapse, acute left ventricular failure or acute renal failure. One treatment session may take hours to complete. It also could present certain other risks: heart failure caused by the rapid transfer of blood, blood loss, acute kidney failure due to temporary major depletion of effective circulating plasma volume, and/or the possible spreading of dangerous diseases such as HIV, hepatitis B, and hepatitis C. The therapeutic method of the present invention avoids these problems. It merely requires an injection, or other equivalent mode of administration, of an antibody composition to the patient.\nHigh dose intravenous immune globulin (IVIG) infusions have also been used in treating certain autoimmune diseases. Previous studies have indicated that IVIG may contain anti-idiotype activity against anti-DNA antibodies, as well as many other autoantibodies (Jordan, S. C., 1989; Silvestris et al., 1994; Mouthon et al., 1996; Silvestris et al., 1996). The effects of IVIG infusions are apparently related to changes in the repertoire of autoantibodies expressed in the patient. This modulation of pathogenic Id antibodies is thought to depend on their specific interaction with the regulatory anti-idiotype molecules that occur naturally in healthy donors. Production of anti-idiotypic antibodies inhibiting the potentially harmful autoimmune repertoire may result from activation of the Id network committed to controlling the secretion of natural autoantibodies by CD5-positive B cells.\nUp until the present time, treatment of SLE with IVIG has provided mixed results, including both resolution of lupus nephritis (Akashi et al., 1990), and in a few instances, exacerbation of proteinuria and kidney damage (Jordan et al., 1989). The cause of this increase is not clear but it is believed that there is increased glomerular deposition of immune-complexed, polyreactive, non-Id-specific IgG antibodies.\nAs can be understood from the above, although there are several treatments for autoimmune disease such as systemic lupus erythematosus, all possess serious disadvantages. Thus, persons afflicted with SLE who show clinical evidence for SLE nephritis need a cost-efficient and safe treatment that will help prevent or ameliorate the tissue damage that leads ultimately to kidney failure and the need for chronic hemodialysis and/or renal transplantation caused by their condition.\nThe subject invention concerns novel compositions and methods for the treatment of antibody-based autoimmune diseases, and in particular, SLE. One aspect of the present invention concerns a therapeutic method for treating patients suffering from, or predisposed to, autoimmune disorders such as SLE nephritis. The method comprises administering an anti-idiotypic antibody composition to a patient afflicted with an autoimmune disease, wherein the anti-idiotypic antibodies selectively immunoreact with autoantibodies bearing the appropriate idiotype, thereby inhibiting the autoantibodies and their destructive autoimmune responses without inducing generalized immunosuppression. The novel anti-idiotypic antibody compositions, prepared in accordance with the procedures of the subject invention, comprise anti-idiotypic antibodies having specificity for pathogenic self-reactive antibodies in a patient\"\"s body, thereby modulating the potential of the self-reactive antibodies to form immune complexes with self antigens and cause harm to normal cells and tissues, particularly within the patient\"\"s kidney filtering systems. One embodiment of the present invention is a method for treating SLE using anti-DNA anti-idiotypic antibody compositions prepared from pooled human intravenous gamma globulin (IVIG). The potential risks and negative side effects associated with other current autoimmune disease therapies are avoided with the present method.\nAnother aspect of the subject invention concerns novel compositions comprising purified human anti-idiotypic antibodies which have binding specificity for self-reactive antibodies that are associated with clinical pathogenesis of certain autoimmune diseases. The antibody compositions are capable of modulating the deleterious effect of the self-reactive autoantibodies on cells and tissues of the affected patient. Specifically exemplified in the present invention is a purified anti-DNA anti-idiotypic antibody composition that can be used to treat patients afflicted with SLE nephritis.\nIn a further aspect, the subject invention is directed toward a method of producing the novel anti-idiotypic antibody compositions which can be used in the therapeutic method of the present invention. The subject antibodies can be produced by adsorbing pooled normal human gamma globulin with a solid phase substrate having antibody molecules attached thereto that selectively bind with anti-idiotypic antibodies present in the human gamma globulin, whereby an antibody/anti-id antibody complex is formed on the solid phase. The anti-idiotypic antibodies from the intravenous gamma globulin preparations are then eluted from the solid phase substrate. The procedure allows for the isolation and enrichment of anti-idiotypic antibodies from the other antibodies present in the pooled gamma globulin preparation."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus for forming an image and outputting an image signal, and a light source unit used with the image forming apparatus.\n2. Related Background Art\nAs one of conventional image forming apparatus, a light source switching type color image sensor is known which reads a color image by applying light beams having three different spectral characteristics to the same color image and outputting image signals. FIGS. 1 to 4 show an example of such an image forming apparatus. This image forming apparatus is constituted by LEDs of red, green, and blue (hereinafter abbreviated as R, G, and B) colors, a short focal point focussing element array, and a sensor array with a plurality of line sensors being disposed in line.\nFIG. 1 is a perspective view showing the image forming apparatus, and FIG. 2 is a cross sectional view of the image forming apparatus. Referring to FIGS. 1 and 2, in the fundamental structure of the image forming apparatus, a light beam 212 outputted from an LED array 211 on an LED substrate 210 mounted on a frame 200 is applied to an original which is in contact with a transparent glass plate 201 mounted on the upper area of the frame 200, and a light beam 213 reflected from the original is applied via an optical system 209 to a sensor array 1 on a sensor substrate 19.\nAs shown in FIG. 3, the LED array 211 has a plurality set of LED chips 211R, 211G, and 211B alternately disposed in line on the LED substrate 210. The LED chips 211R, 211G, and 211B emit RGB light beams, and light of each color of RGB can be independently turned on and off. The optical system uses a short focal point focussing element array, for example, a product name \"Selfoc Lens Array\" manufactured by Nippon Sheet Glass Co., Ltd.\nAs shown in FIG. 4, the sensor array 1 has a plurality of line sensors 2-1, 2-2, . . . , 2-15 disposed in line on the sensor substrate 19, the line sensors being covered with a protection film 206. A tight contact type multi-chip image sensor fundamentally reads an image by applying a light beam reflected from an original to a sensor array and focussing an image of the same size as an original. Therefore, the length of the sensor array 1 is required to be equal to or longer than the width of an original.\nThe length of the sensor array 1 changes with the size of an original to be read, and the number of line sensors of the sensor array 1 changes. For example, for reading of an A3 size original, the sensor array has fifteen line sensors assuming that the length of each line sensor is 20 mm.\nThe sensor substrate 19 is coupled via a flexible substrate 208 to another substrate 203 on which a connector 202 for input/output of a power source and control signals is mounted. The substrate 203 is fixedly mounted on the frame 200 by means of screws 207.\nNext, the read operation of the image forming apparatus will be described. First, data for correcting shading error is read, the shading error being generated by a variation of line sensor sensitivities and a variation of emission of a light source. In reading shading correction data, LED chips 211R, 211G, and 211B are sequentially turned on to read a white reference plate built in the image forming apparatus, and the output signals of the image sensor are temporarily stored in memories provided for respective colors.\nBy reading the sensor output signals r1 for LED 211R, g1 for LED 211G, and b1 for LED 211G obtained by independent emission of RGB light sources and stored in the memories, the gain of each color is adjusted to satisfy the condition of r=g=b where r, g and b are sensor output signals for RGB colors obtained when the white reference plate is again read.\nIn reading an original with a light source switching color image sensor, it is necessary to independently apply RGB light beams to the original in order to obtain three RGB signals. To this end, a frame sequential method and a line sequential method are used. With the frame sequential method, LEDs of one color among RGB colors are turned on to sub-scan the whole frame of the original, this operation being repeated for the other two colors. With the line sequential method, LEDs of three colors are sequentially turned on for each line of an original to sub-scan the whole frame. Both the methods can obtain RGB signals of the whole area of an original to reproduce a color image.\nThe ideal spectral characteristics of RGB light sources for a light source switching color image sensor will next be described. A G light source is used by way of example. As shown in FIG. 5, it is assumed that an original image is read by using three G light sources each having a different spectral characteristic. A light source G6 does not contain light in the wavelength ranges from near 480 to near 500 nm and from near 570 to near 590 nm, as compared to the light source G7.\nTherefore, if colors a and b shown in FIG. 6 having different spectral characteristics only in the wavelength range near 500 nm are read by using the light source G6, a difference of the spectral characteristics between the colors a and b cannot be discriminated and generally a same G signal is obtained for the colors a and b.\nIf the B light source having a shorter wavelength than the G light source does not contain light in the same wavelength range as the light source G6, then the colors a and b cannot be discriminated. In order to improve color discrimination between various colors contained in a color original, the spectral characteristics of the RGB light sources are required to cover the whole visible light range.\nNext, a difference of color reproduction between the light sources G7 and G8 will be described. Light of the light sources G7 and G8 covers the same wavelength range, and only the energy distribution in the wavelength range is different. Color spaces of a light source switching color image sensor using the light sources G7 and G8 are shown in FIG. 7.\nThe diagram shown in FIG. 7 is called a CIE-xy chromaticity diagram. In FIG. 7, all colors are contained in an area surrounded by a solid curve line representative of a spectrum locus or reddish-purple line. Triangles in this area represent color spaces of the color image sensor. An output G.sub.OUT of an image sensor when an original is applied with light from the light source G7 or G8 is given by the following equation. EQU G.sub.OUT =.intg.G7(.lambda.)(or G8(.lambda.))S(.lambda.)d.lambda.\nwhere G7(.lambda.) represents a spectral emission characteristic of LED G7, G8(.lambda.) represents a spectral emission characteristic of LED G8, and S(.lambda.) represents a spectral sensitivity characteristic of a line sensor.\nColor reproduction is made not by measuring the detailed spectral reflection characteristic of an original, but by using RGB signals. As seen from FIG. 7, the color image sensor using the light source G7 has a broader color space than the light source G8. The spectral characteristic of the light source G7 for a light source switching color image sensor is more desirable than the light source G8.\nThe ideal spectral characteristics of RGB light sources for a light source switching color image sensor are therefore required to have as broad color spaces as possible and cover the whole wavelength range. An LED light source has many advantages such as compact size, high response speed, and good reliability, over other tubular type light sources. Therefore, it is suitable for use with a light source switching color image sensor.\nColor reproduction using RGB signals of an LED light source switching color image sensor is, however, associated with some problems. FIG. 8 shows an example of a color space of a conventional LED light source switching color image sensor. As seen from FIG. 8, the color space of the image sensor is rather narrow as compared to various colors in a natural world. This results from that the spectral characteristic of an LED of G color is positioned too near the long wavelength side and that there is a wavelength range having too small light emission. There are LEDs for three colors having the spectral characteristics which can solve the above problems and realize ideal color reproduction. However, such LEDs are very expensive and the manufacturing cost of an image forming apparatus becomes too high. In contrast, general LEDs used for display devices or other devices are mass-produced and relatively cheap. If such display LEDs are used as light sources of an image sensor, the cost can be reduced considerably.\nHowever, the display LED has a sharp spectral characteristic with a small full width at half-maximum. It is necessary for a light source of a light source switching color image sensor to cover the whole visible light range by using LEDs of three RGB colors. Therefore, if display LEDs are used as the light sources for an image sensor, they are associated with the above problems such as a wavelength range with an extremely small emission amount and poor color reproduction, because the display LEDs have too narrow full width at half-maximum.\nIn a conventional image forming apparatus, LED chips of three RGB colors are disposed at an equal pitch. Therefore, as shown in FIG. 9, an incident angle of light from an LED for each color is different for each color at an arbitrary point on an original. As a result, the optical information of an original supplied to a sensor pixel train, i.e., the intensity of vertical components of a reflected light beam, is different at each point on the original. From these reasons, even if an original having a uniform density is read, the color component ratio is different at each point of the original to thereby result in color shade. As shown in FIG. 10, the size of shadow at the corner of a convex portion of an uneven original such as an original with a pasted sheet changes with the color component, coloring the shadow at the corner."}
{"text": "The invention relates to a weft stop motion, or detector, for looms in which the sensor element responds or is sensitive to an electric charging of the weft thread without contact and also to looms in which the weft break stop motion of the present invention is used.\nA concept for thread detection is known from German Patent Specification 3,758,403, for example. Various embodiments of electrostatic transformers are disclosed therein. These sensors are mainly used in air-jet looms. The weft thread is electrically charged during its removal from the weft thread supply because of the resultant friction and also during the weft insertion because of friction with the air. The electrostatic detection registers the presence of a textile fiber which is moving past and is electrically charged in this way, and the passage of the tip of an inserted weft thread in particular can also be detected. Weft break stop motions are used in the weft channel of the loom. The known embodiments are relatively large and heavy and are frequently constructed in the form of a confusor drop wire. The high-speed air-jet looms having a correspondingly high beat-up speed of the reed which are commonly used nowadays produce high vibration and acceleration loads on the known weft break stop motions, so that the known embodiments are no longer suitable for use on air-jet looms or the resultant electrical signals are very noisy."}
{"text": "1. Field of the Invention\nThis invention relates to molding compositions and forming processes for normally rust-prone iron-based metal alloy powders, and articles produced therefrom. Metal alloy systems that can be successfully formed using the processes of the invention, include elemental iron and iron alloys, including low and medium alloy steels, tool steels, and a number of specialty iron-base alloys.\n2. Description of the Related Art\nA widely used process for forming metal powders into complex three dimensional shapes is Metal Injection Molding (MIM). The steps of fabrication of metal or ceramic-metallic (CERMET) parts are the following:                i. Metal and/or ceramic powders are blended with a thermoplastic binder material to create an injection molding feedstock with thermoplastic properties.        ii. The thermoplastic feedstock is injection molded in a fluid state using methods and tools typical of conventional plastic injection molding, and removed from the mold in a solid state.        iii. The “green” state as-molded parts are subjected to thermal and/or chemical processes to remove the binder phase.        iv. The resulting “brown” state metal or CERMET parts are sintered at higher temperatures to effect consolidation and densification of the molded object.        \nSeveral methods, processes, and binder systems have previously been described for fabrication of rust prone iron-based metal alloys and CERMET materials containing them. Each of these processes has one or more disadvantages that prevent important applications.\nFor example, commonly utilized polymer or wax binder MIM processes, such as the methods described by Achikita et al. in U.S. Pat. No. 5,250,254, while they work well with rusting iron alloys, are limited to small parts, weighing no more than a few hundred grams, and with maximum section thickness of less than 10 millimeters. These limitations are imposed by the difficulties associated with binder removal prior to sintering. The manufacture of larger parts is prevented or rendered uneconomical by dimensional instability, cracking, or simply the long times needed for binder removal from larger and thicker sections. In addition, great care must by taken when using wax or resin binders to avoid an undesirable out-of-specification increase in the carbon content of the alloy as a result of incomplete removal of the hydrocarbon binder phase.\nFanelli et al., in U.S. Pat. No. 4,734,237, disclose agaroid-based aqueous binders for molding of metal and ceramic powders. The development of aqueous-binder molding compositions, including those disclosed by Fanelli et al., has largely removed the part size restrictions imposed by wax and polymer binders, since the binder phase in these largely consists of water which is easily removed by evaporation under ambient conditions. In the special case of agar-based binders, the carbon content problem associated with wax and polymer binders is also reduced since the agar component of the binder is largely gasified at relatively low temperatures during the early stages of the sintering cycle. Further reduction in carbon content is easily achieved by employing an oxidizing atmosphere in the early stages of the sintering heat treatment as taught by Zedalis in U.S. Pat. No. 5,985,208. Carbon content can also be reduced by heat treatment in hydrogen as taught by Wu et al., “Effects of residual carbon content on sintering shrinkage, microstructure and mechanical properties of injection molded 17-4 PH stainless steel,: Journal of Materials Science, 37 (2002) pp. 3573–3583.\nZedalis et al., in U.S. Pat. No. 6,268,412, incorporated herein by reference to the extent not incompatible herewith, disclose molding compositions and processing steps for injection molding of non-rust-prone stainless steel articles using water-base agaroid binder systems. Stainless steels, a family of iron-based alloys containing between 10.5 and 28 atomic % chromium, are compatible with water-based binder systems, since the high chromium content confers great resistance to oxidation in the presence of water.\nWhen rust-prone iron-base alloy powders are substituted for the stainless steel powders in the process taught by Zedalis, the resulting molding feedstock is chemically unstable and must be molded and dried within hours, or the water will react with the iron-base alloy powder to form rust, thereby substantially altering and degrading the Theological properties, as-molded strength, sintering, and shrinkage behavior of the feedstock.\nIt is commonly observed that ferrous alloys progressively oxidize or rust in the presence of air and moisture. The essential chemistry of rust formation, as described in The Metals Handbook, Volume 1, 8th Edition, published by the American Society for Metals (1961) p257, follows. In the first step of the reaction, iron reacts with water to form ferrous and hydroxyl ions and hydrogen:Fe+2H2O═Fe+++2OH−+H2  (1)\nIn a second step, oxygen, if present, reacts with the ferrous ions to produce ferric ions which precipitate out of solution as insoluble ferric hydroxide FeO(OH), otherwise known as rust. Since the rust deposit does not form a protective layer, reaction 1 is free to proceed until the metallic iron is consumed or equilibrium is reached.\nThe equilibrium constant for reaction 1 is:K=[Fe++][OH−]2PH2  (2)where the square brackets indicate the concentration of the species and PH2 is the partial pressure of hydrogen.\nEquation 2 suggests that the equilibrium concentration of Fe++ can be suppressed by increasing the hydroxyl ion concentration, equivalent to increasing the pH, and/or increasing the hydrogen partial pressure.\nRusting can be further inhibited by passivation of the exposed ferrous alloy surface. Typically, passivation involves a thin but impervious layer of iron oxide formed, in-situ, by reaction of the iron with oxidizing ions. Pourbaix, in Atlas of Electrochemical Equilibria in Aqueous Solutions, Pergamon Press, New York (1966) p. 312 states that passivation of iron is difficult at a pH below 8, relatively easy at a pH above 8 and very easy between pH 10 and 12. Above pH 13, according to Pourbaix, iron will corrode by hyperferrate ion formation. Passivation of ferrous alloy surfaces is typically effected by the addition of oxidizers to aqueous environments. For example, nitrite and nitrate salts have been used in this manner as rust-inhibiting additives in cooling water and other process water applications. pH buffers, salt solutions formed by reaction of weak acids with strong bases, are frequently employed with nitrites and nitrates to maintain pH in the proper range. The Metals Handbook, Vol. 1, 8th Ed., American Society for Metals, P. 279, 1961 states that sodium nitrate-borate combinations have been used to inhibit corrosion in diesel engine cooling systems and in low pressure, hot water recirculating systems. In this case, sodium borate, a salt formed by reaction of the weak acid H3BO3 with the strong base NaOH, supposedly functions as a pH buffer. In a similar fashion, calcium nitrite is frequently added to concrete formulations to inhibit rusting of embedded steel reinforcing bars. In this case, the desired alkaline environment is synergistically provided by the calcium oxide component of the Portland cement concrete.\nInterestingly, various metal borate additives to enhance the gel strength and viscosity of polysacharide-based aqueous binders for molding have been disclosed by Sekido et al. in U.S. Pat. No. 5,258,155, and Fanelli et al. in U.S. Pat. No. 5,746,957. Anions of boric acid, acting in concert with metal cations are thought to induce crosslinking of the agar polysaccharide molecules, thereby substantially increasing the viscosity of the agar-water sol and the strength of the gel. Fanelli et al. teaches that calcium borate, magnesium borate, zinc borate, ammonium borate, tetraethyl ammonium borate, tetramethyl ammonium borate, and boric acid are preferred gel strengthening additives for agaroid binder powder injection molding of ceramic and/or metal powders.\nSekido et al., broadly teach the use of sodium borate for gel strengthening but, Sekido does not define sodium borate in useful chemical terms. That is, neither the preferred concentration range nor the preferred stoichiometry range (i.e., the preferred atomic fractions of sodium and boron) are specified.\nFor the purposes of the present invention, it is important to clearly distinguish between the different borate salts of sodium and potassium. According to the CRC Handbook of Chemistry and Physics 56th edition, CRC Press, Cleveland Ohio (1974), two crystalline sodium borate salts are known, sodium tetraborate (Na2B4O7) and sodium metaborate (NaBO2). Moreover, both of these may occur as anhydrous or hydrated salts. The most familiar sodium borate salt is the mineral borax, or sodium tetraborate decahydrate (Na2B4O7.10H2O). In aqueous sodium borate solutions, of course, one is not limited to these fixed stoichiometry compounds and a continuous range of boron to sodium ratios can be obtained between the endpoints NaOH and H3BO3. Similarly, potassium tetraborate, potassium tetraborate tetrahydrate, potassium metaborate and other fixed stoichiometry crystalline potassium borate compounds are known, but any boron to potassium ratio can be obtained in solution. The various sodium and potassium borate salts are formed by reaction of the weak acid H3BO3 with the strong bases NaOH and KOH. For example:H3BO3+NaOH═NaBO2+2H2O  (3)\nFor clarity in describing various borate salts herein, the molar ratio of H3BO3 to NaOH will be used to specify the stoichiometry of sodium borate salt solutions, and the molar ratio of H3BO3 to KOH will be used to specify the stoichiometry of potassium borate salt solutions. These ratios are the same as the atomic ratios of boron to sodium and boron to potassium. Thus, Na2B4O7 has a B:Na ratio of 2:1 while NaBO2 has a B:Na ratio of 1:1. We will also at times use the mole fraction of H3BO3 used to make the salt solution, defined as (moles H3BO3)/(moles H3BO3+moles (Na,K)OH). Thus, a solution of NaBO2 has a mole fraction of H3BO3 of 0.5 or 50%, and a solution of Na2B4O7 has a mole fraction of H3BO3 of 0.66 or 66%. These conventions are convenient for the synthesis of various sodium and potassium borate salt solutions from boric acid (H3BO3), which is available as a crystalline solid, and the respective sodium and potassium hydroxides, which are readily available as solutions of specified molar concentration.\nThus, the concentration and stoichiometry of a solution of any sodium or potassium borate salt can be fully described by specifying the equivalent molar concentrations of H3BO3 and NaOH or KOH in the solution.\nFor example, Sekido, in his Example 1, used a combination of agar and an aqueous solution of sodium borate as a binder for 316 stainless steel powder. The concentration of sodium borate in the water was about 0.3 wt. %. Presumably the sodium borate used was common borax (sodium tetraborate decahydrate). The molar concentration of Na2B4O7.10H2O was therefore 0.0079 moles/liter, the equivalent molar concentration of H3BO3 was four times this or 0.0316, and the equivalent molar concentration of NaOH was twice that of Na2B4O7.10H2O or 0.0158.\nBehi et al. in U.S. Pat. No. 6,261,336, specifically addressed the problem of rust formation in aqueous agar binder injection molding feedstocks containing rust-prone ferrous alloy powders, and taught that these materials can be stabilized against rust formation by the addition of alkaline sodium silicate to the aqueous binder. It was shown by Behi that carbonyl iron powder feedstocks containing appropriate amounts of sodium silicate are somewhat stable against rust formation and attendant hydrogen evolution, and that the stability is further enhanced by the addition of potassium borate. The sodium silicate is thought to function by reacting with the iron surface to form a barrier layer of iron silicate and the potassium borate in this application apparently serves as a pH buffer similar to the use of the sodium borate/nitrite combination discussed above. Behi cites potassium tetraborate and potassium tetraborate tetrahydrate as preferred potassium borate compounds and gives a preferred borate concentration range of from about 0.01 to about 0.2 weight % of the composition (which would correspond to about 0.125–2.5% weight % relative to the aqueous solvent at a typical moisture content of 8 wt. %). While Behi's sodium silicate/potassium borate stabilized feedstocks certainly represent an improvement over unstabilized iron-based aqueous binder feedstocks, experience with the sodium silicate stabilized feedstocks has revealed that the long term chemical stability is marginal, and that the sodium silicate addition renders the feedstock pellets somewhat tacky and difficult to feed through the hopper of an injection molding machine. Moreover, residual SiO2 and/or iron silicate inclusions, resulting from decomposition of higher loadings of the sodium silicate during sintering, may be undesirable for applications requiring maximum ductility and fatigue resistance in the final sintered steel part.\nMore recently, Morris, in U.S. Pat. No. 6,689,184, has disclosed stabilization of rusting iron aqueous molding feedstocks using a combination of borate and nitrate/nitrite salts. One disadvantage of this system is that the nitrate and nitrite salts serve as nutrients for a range of micro-organisms. Another disadvantage is that the nitrate and nitrite salts may tend to oxidize minor alloy components such as silicon and chromium during the elevated temperature sintering process.\nThus, a need remains for new materials and methods enabling molding of rust prone iron-based alloys that avoid the size limitations of the prior art wax and polymer based binders, and the processing and ductility limitations of sodium silicate and nitrite/nitrate stabilized aqueous binders."}
{"text": "1. Field of the Invention\nThe present invention relates to a portable wireless terminal. More particularly, the present invention relates to an antenna device equipped in a portable wireless terminal.\n2. Description of the Related Art\nAs the electronic communication industry develops, portable wireless terminals are becoming lighter, slimmer, smaller, and more multi-functional. For example, a speaker device capable of realizing melodies of various harmonies is installed, and a color display device with millions of pixels is implemented. Also, in addition to a call function, a portable wireless terminal now typically provides a music listening function through a Moving Picture Experts Group (MPEG) Audio Layer 3 (MP3) Player (MP3P). Furthermore, the portable wireless terminal provides not only various game contents using the display device but also a function of receiving a radio signal, a Digital Multimedia Broadcasting (DMB) signal, etc.\nIn general, a portable wireless terminal uses an antenna that transmits and receives signals. If it is desired to provide the portable wireless terminal with greater radio wave transmission/reception, the antenna is designed to be large so as to have high directivity. In such a case, the spacing distances between other parts within the portable wireless terminal decrease in order to maintain the lightness and smallness of the portable wireless terminal. Therefore, it is difficult to secure the space for the larger antenna and improve the transmission/reception performance of the terminal. In addition, a Transmit (Tx) output of a designed antenna should not be allowed to exceed a reference value permitted in each country.\nFor example, an antenna for transmitting and receiving a signal in a Frequency Modulation (FM) frequency band may be implemented with a single antenna (e.g., an earphone connected to an external connector). In this case, a problem arises in that a Transmit (Tx) output of the antenna for transmitting and receiving the signal in the FM frequency band exceeds a reference value permitted in each country. FIG. 8 is a graph illustrating an antenna performance of a conventional antenna transmitting and receiving singles in the FM spectrum. FIG. 9 is a table illustrating an FM Tx output regulation specified by the USA and Europe. In a case where an earphone antenna is used, the Tx output is 98 dB which exceeds an output regulation specified by each country as shown in FIG. 9. Accordingly, a need exists for an antenna device that increases the reliability of communication while not exceeding output regulations of each country."}
{"text": "This invention relates to data communication systems for interconnecting an end user machine with a remote server (e.g., an Internet server) for the two-way transmission of data packets. More particularly, the invention relates to wireless communication links, such as cellular packet networks, in which mobile subscriber units may be switched between base stations.\nA communication system of this type typically transports a sequence of data packets over a TCP connection or the like between an end user machine coupled to the subscriber unit and a server (e.g., an Internet server) coupled to the base stations through a fixed network. In the wireless portion of such system, data packets from the server flow to the subscriber unit through the base station that registers the strongest signal strength as measured, e.g., by a beacon or pilot signal received by the subscriber unit. If a subscriber unit that is initially serviced by a first base station roams through an area where the signal strength is stronger from a second base station, the subscriber unit typically requests a change of transmission path (e.g., a “handoff”) from the first base station to the second station.\nPropagation delays, data bit errors and the like are normal on wireless communication links. Such phenomena can cause loss or delay of acknowledgment signals that are successively generated by the end user machine in response to successive bytes contained in data packets received by the end user machine from the server. Each acknowledgment signal contains a first identifying portion indicative of the corresponding byte received by the end user machine and a second portion advertising the then-current size of the receive window of the end user machine.\nThe loss or delay of acknowledgment signals is often interpreted as congestion on the network by the applicable TCP protocols which were designed primarily for end-to-end wired networks. As a result, the server may be switched into a so-called congestion avoidance or slow-start mode, which can drastically reduce throughput of data packets on the system even when no congestion is present.\nWhile known techniques involving, e.g., modification of the network protocols, attempt to mitigate the effects of such loss of throughput in wireless systems, they frequently add complexity such as the splitting of the TCP connection between the end user machine and the server. More importantly, the effectiveness of such techniques is greatly diminished during periods of handoff."}
{"text": "The prior art is well documented with examples of extraction grip or retrieval devices. The purposes of such devices include their use and application for facilitating the removal of an ammunition holding magazine, such as which is secured to a projectile firing device.\nA first example is shown in U.S. Pat. Nos. 6,748,689 and 6,883,261, both to Fitzpatrick, and which discloses integral extensions for aiding in the extraction of ammunition magazines from ammunition pouches. The sides of the ammunition magazine are extended, either by molding or affixing a handle directly to the sides of the magazine, to thereby provide a grasping handle.\nOther examples include U.S. Pat. Nos. 6,634,131 and 6,212,815, also to Fitzpatrick, and which describes another type of magazine extraction grip including a sleeve of resilient material molded in the general shape of a magazine, however exhibiting a smaller inner circumference than that associated with the magazine to require the band to be stretched over the magazine. The top of the band exhibits a smaller inner circumference than the lower part of the band and extending from the top is a handle for permitting a users finger to wrap there-around to extract the magazine.\nYet further examples of integral magazine extraction extensions are shown in U.S. Pat. Nos. 7,207,131, 7,174,666 and 6,481,136, also to Fitzpatrick. A handle is permanently attached to an existing or modified floor plate of a magazine. In one example, a substitute floor plate is provided with a molded projecting handle. In another, a handle with a grip and a terminal end is provided with an attachment structure thereupon. Existing floor plates can also modified by cutting anchoring holes to allow for attachment of the handles, and without hindering use in an ammunition magazine or molded with the anchoring holes."}
{"text": "The present invention relates to handrests or maulsticks used by painters for resting a painter\"\"s hand while painting. More particularly, the present invention relates to an adjustable handrest which is capable of removable attachment to a canvas frame to enable a painter to rest his or her brush hand while painting a design or picture on a canvas stretched on the canvas frame.\nA common working medium for artists, sign painters, cartographers, delineators and the like is a sheet of canvas, cardboard or work surface of other material which is mounted on a hollow rectangular frame. The frame is supported in a generally vertical plane on an easel or other support to provide free and flexible access of the artist to the work surface and to orient the work surface in a position for viewing. The artist typically applies the paint or other medium by freehand to the work surface using a brush.\nThroughout the course of preparing a painting or other work of art on canvas, it is often necessary to apply various colors contiguously to previously-applied, but still wet, colors. This, as well as the occasional requirement of drawing straight lines on the canvas, requires the artist to maintain great steadiness of hand while applying the medium to the canvas. Furthermore, freedom of movement of the artist\"\"s hand, which is facilitated by steadiness while applying the medium to the canvas, provides to the artist full expression of his or her theme.\nVarious devices are known in the art for assisting an artist in steadying his or her hand while applying a medium to a work surface. Patents of interest in this regard include U.S. Pat. Nos. 289,700; 518,761; 1,422,641; 2,496,276; 2,814,142; 3,101,568; 3,815,856; 4,188,006; 3,972,133; 4,088,290; 4,685,644; 5,141,198; 5,172,883; 5,193,772; 5,299,772; and 5,765,791.\nThe present invention includes an adjustable handrest for artists which is capable of attachment to a frame for supporting a worksurface on which a painting or other work of art may be applied. The handrest includes an elongated hand support which spans the work surface and each end of which is independently vertically adjustable with respect to the other. The handrest is equally well-adapted for use by left-or right-handed users."}
{"text": "This invention relates to a method for building a wall or a part thereof in the ground, a system for carrying out the method, and a wall building device for use in the method and system.\nKnown tunnel wall building devices are e.g. described in the xe2x80x9cHandbook of Mining and Tunnel Machineryxe2x80x9d, Barbara Stack, 1982, published by John Wiley and Sons, pages 415-417. These known tunnel wall building devices comprise a cutting face at the front of a cylindrical shield. The tunnel wall building device is pushed in the direction of the advance of the shield by hydraulic cylinders acting between the rear side of the shield and the constructed tunnel wall.\nA first drawback of the use of such tunnel wall building devices is the criticality of the control over the stabilization of the ground in front of the cutting face. Too much excavation will cause local collapse of the ground, while too great a push force will cause undesirable settlements of the ground.\nA second drawback is the required axial support of the tunnel wall building device on the tunnel wall. In case of the use of prefabricated elements the axial support force might be a governing load case.\nIn other methods the tunnel wall is made of extruded concrete. In this respect reference is made here to EP-A-0 354 335 describing axially supporting a boring shield by the extruded concrete tunnel wall through formwork elements used for forming the inner side of the tunnel wall. The tunnel wall building method according to this publication is discontinuous, taking away formwork elements where the concrete has hardened, and adding the formwork elements, after cleaning thereof, directly behind the boring shield. This method is expensive and time-consuming as a result of the use of the formwork elements, which need to be handled behind the boring shield, and can only be removed when the concrete they are supporting has hardened sufficiently, which may take considerable time.\nEP-A-0 483 445 describes a continuous tunnel wall building method, extruding concrete to form a tunnel wall, and using a sliding formwork arrangement requiring additional means to transfer the axial boring forces to the hardened part of the tunnel wall and to facilitate continuity.\nIn all previous methods the management of the ground water level during the fabrication of the tunnel wall may present serious problems.\nThe object of the present invention is to provide a method, system, and device for building a wall or a part thereof in general, and a tunnel wall in particular which avoid the risk of collapse or undesirable settlement of the ground in front of the wall building area.\nAnother object of the present invention is to provide a wall building method, system, and device which do not rely on the wall for providing an axial support for the wall building activities.\nYet another object of the present invention is to provide a wall building method, system, and device which can be continuous and do not use formwork elements.\nA further object of the present invention is to provide a wall building method, system, and device which are virtually unaffected by the ground water level.\nTo reach the above objects, in the method according to the invention a wall building device having cross-sectional dimensions which are substantially equal to the dimensions of at least a part of the cross-section of the wall is pulled through the ground, at least a part of the wall being formed by injecting a hardenable material behind the wall building device. In the process of making the wall, the wall is supported on all sides by the ground surrounding the wall (and, in the case of a tunnel wall, by the ground filling the tunnel wall). Consequently, the wall requires no direct strength, and can e.g. be made from fiber concrete with a normal hardening time. The excavation of the ground at one side of the wall can be done after the completion of the tunnel wall, and the stability of the excavation presents no problem at all. Different shapes of wall cross-sections are possible, but not limited to: (semi-)circular, elliptical, rectangular, triangular, flat. Since the wall building device is advanced by pulling, no reaction force is applied on the constructed wall. The advance of the wall building device is therefore independent from the structural strength of the wall at any given moment.\nIn a preferred embodiment of the wall building method according to the invention, the following steps are performed: drilling one or more holes in the ground, on or parallel to the projected path of the wall, each of the one or more holes being drilled by means of a drill string; connecting one end of each drill string at the end of the drilling operation to the wall building device; and pulling the wall building device through the ground by retracting the one or more drill strings. Alternatively, the wall building method may comprise the steps of: drilling one or more holes in the ground, on or parallel to the projected path of the wall, each of the one or more holes being drilled by means of a drill string; connecting one end of each drill string at the end of the drilling operation to one end of a casing string; pulling each casing string through the ground by retracting the drill string connected thereto; disconnecting each casing string from the corresponding drill string; connecting one end of each casing string to the wall building device; and pulling the wall building device through the ground by retracting the one or more casing strings. Accordingly, depending on the kind of wall to be built, in particular its cross-sectional area and its length, one or more drill strings and or one or more casing strings are used to pull the wall building device through the ground, using standard directional drilling techniques to bring the drill and casing strings into the ground. The directional drilling techniques are capable of very accurately following the projected path of the wall. The accuracy of the paths of the casing strings can be further improved by using a single drill string and by means of a connecting assembly connecting the ends of several casing strings to one end of the drill string and transversely spaced therefrom, which drill string is then pulled through the ground at the other end thereof.\nIn a preferred embodiment each casing string comprises an inner string and an outer string enclosing the inner string, the inner string being adapted to transfer the pulling force required by the wall building device, and the outer string being adapted to provide low friction forces when moving the inner string relative to the outer string. The outer string serves as a guide for the inner string when pulling the wall building device through the ground by the inner string. The friction forces between the inner string and the outer string can be still further reduced by supplying a fluid to the space between the inner string and the outer string.\nPreferably, the inner string of the casing string is made of steel providing the strength needed, while the outer string is made of plastics, e.g. polyethylene providing an excellent low cost and low friction separation wall between the inner string and the ground.\nWhen moving the wall building device through the ground, a hardenable material is injected behind it to form the wall. The wall material preferably is supplied to the wall building device through at least one of the one or more drill strings/casing strings used for pulling the wall building device through the ground. Additionally, at least one of the one or more drill strings/casing strings may contain at least one line for supplying energy, at least one line for controlling and/or monitoring the wall building device, at least one duct for supplying a drilling fluid to the wall building device, and/or at least one duct for discharging ground removed by the wall building device. Alternatively, or in addition thereto, the wall material and the drilling fluid may be supplied to the wall building device through ducts extruded in the wall. Such ducts can also be used for accommodating lines for supplying energy to the wall building device, or for controlling and/or monitoring the wall building device. Further, such ducts can be used for discharging ground removed by the wall building device.\nFor an improved positioning accuracy of different wall parts relative to each other, preferably the wall building device is adapted to build a wall or wall part provided with wall guide means, the or a further wall building device being provided with means for engaging the wall guide means. In this way the latter wall building device is allowed to exactly follow the path of the wall or wall part already in place in the ground. The wall guide means preferably comprise a guide slit, which may have an essentially L-shaped cross-section. More generally, in a preferred embodiment at least one guide member is provided in the ground, the wall building device is adapted to be guided along the guide member through the ground.\nWhen a wall building device is used in a multi-pass mode, i.e. for subsequently passing through the ground along adjacent paths for building a wall from a plurality of separately built wall sections, preferably the wall building device is adapted to be connected to the leading end of a pulling string for pulling the pulling string into the ground simultaneously with pulling the wall building device through the ground. The pulling string is used for pulling the or a further wall building device through the ground in a next passage. For this purpose, the wall building device preferably has an essentially Z-shaped cross-section.\nAn essentially flat wall can be built in the ground by at least two wall building devices which at at least one side thereof are adapted to be coupled to another, which wall building devices further are adapted to be pulled through the ground in a forward direction, and in a direction at an angle to the forward direction.\nFor improving the quality of the wall, the wall material preferably is injected between one or more membranes defining one or more sides of the wall or wall part formed by the wall building device. The membrane or membranes can be stored in the wall building device in a folded or rolled-up form, and be unfolded or unrolled when injecting the wall material.\nThe friction between the wall building device and the surrounding ground, and the resistance the ground offers to the wall building device when pulling the wall building device through the ground is preferably reduced by vibrating, lubrication, jetting and/or removing the ground area adjacent to the front part of the wall building device.\nThe claims and advantages will be more readily appreciated as the same becomes better understood by reference to the following detailed description and considered in connection with the accompanying drawings in which like reference symbols relate to the same parts or parts having the same function. Arrows without reference numerals indicate normal directions of movement."}
{"text": "Augmented reality (AR) technology augments an image of a real-world environment (that is, reality) by superimposing supplemental information (such as, pictures, videos, three-dimensional (3D) models, and other sensory enhancements) onto the real-word environment. The AR technology overlays virtual objects onto the image of the real world, enhancing a user's perception of reality and providing a user with an immersive, interactive experience."}
{"text": "An integrated circuit comprises a substrate on the surface of which and/or inside of which are produced electronic components to form one or more electronic chips.\nThese integrated circuits are produced by using collective fabrication methods often called “microelectronics” methods. For example, these methods implement a machining of the substrate or layer by photolithography and etching (for example DRIE (Deep Reactive Ion Etching)) and/or a structuring by epitaxial growth and deposition of conductive material. By virtue of these microelectronics methods, the fabricated electronic components are small. In many cases, these components have dimensions of a micrometric or nanometric order. The dimension of micrometric or nanometric order is generally less than 10 μm and, typically, less than 1 μm.\nThe substrate can be homogeneous, such as a block. Or, the substrate can be heterogeneous, such as a plurality of stacked thin layers. In this case, the electronic components are produced on the surface and/or inside this substrate and the electronic chip or chips are arranged at least partly on the surface of this substrate.\nThe substrate can also be a “multilayer” substrate. Such a substrate is made of up of a stack of substrates assembled one on top of the other. The stacked substrates can themselves be homogeneous or heterogeneous. Here, to distinguish this substrate from the stacked substrates of which it is composed, the stacked substrates are called “layers.” In the case of a multilayer substrate, the stacked layers are not thin layers. In particular, each stacked layer exhibits a sufficient rigidity in itself to be handled without being bonded onto another substrate. To this end, typically, the stacked layers have a thickness at least 10 or 100 times greater than the thickness of the thin layers. For example, the thin layers have a thickness of less than 1 μm or 0.1 μm, whereas the stacked layers have a thickness greater than 5 or 10 or 100 μm. Several of these layers can include electronic components on the level of the assembly interface between two successive layers of the substrate. These electronic components produced on the surface of the stacked layers constitute one or more electronic chips. The integrated circuit then comprises a stack of electronic chips, stacked one on top of the other. These integrated circuits in which a plurality of electronic chips are stacked one on top of the other are often called “3D integrated circuits.”\nOne of the main limitations on performance of integrated circuits is the quantity of heat that they produce and which has to be removed. This problem is even greater in 3D integrated circuits.\nIn order to evaluate, as correctly as possible, the heat from an integrated circuit, it is useful to know the places inside this integrated circuit where the heat is concentrated and reaches a maximum. These places are called “hot spots.” It is also important to know the diffusion from these hot spots inside the circuit.\nOne of the difficulties commonly encountered is that of knowing both the positions of these hot spots and the diffusion from these hot spots. This difficulty is exacerbated when the hots spots are inside the substrate, especially if they are more than 5 μm away from an outer face of the substrate.\nIt is known to measure temperature on the surface of the substrate, but now how to directly measure temperature inside the substrate. Thus, currently, to estimate the temperature of a hot spot situated inside the substrate, simulation software is used. Such software is based on mathematical models that make it possible to predict the flow of the heat fluxes inside the substrate. To improve the estimation of this software, the latter can also use, as input data, the temperatures measured on the surface of the substrate. However, these models can prove imprecise, notably in estimating the dynamic aspects of certain thermal phenomena.\nThe prior art is also known from: FR2878077A1, DE102010029526A1 and US2010/219525A1."}
{"text": "1. Field of the Invention\nThis invention relates generally to a heat-assisted magnetic recording (HAMR) type of magnetic recording disk drive, and more particularly to a HAMR disk drive with a sensor for measuring laser power.\n2. Description of the Related Art\nHeat-assisted magnetic recording (HAMR), sometimes also called thermally-assisted recording (TAR), has been proposed. In HAMR disk drives, an optical waveguide with a near-field transducer (NFT) directs heat from a radiation source, such as a laser, to heat localized regions of the magnetic recording layer on the disk. The radiation heats the magnetic material locally to near or above its Curie temperature to lower the coercivity enough for writing to occur by the write head. HAMR disk drives have been proposed for conventional magnetic recording disks where the magnetic recording layer on the disk is a continuous layer of magnetic recording material. HAMR disk drives have also been proposed for bit-patterned media (BPM) where the magnetic recording layer is patterned into small isolated data islands, each island containing a single magnetic “bit” and separated from neighboring islands by nonmagnetic spaces.\nIt is important during writing that the output power at the NFT be within a predetermined range. If the laser power setting is too low and thus the NFT output power too low, the desired data bit will not be magnetized. If the laser power setting is too high and thus the NFT output power too high, bits adjacent to the desired data bit will also be magnetized. While a photo-detector can be used to monitor the laser power, the additional cost makes this an unattractive solution. An electrically conductive thermal sensor has been proposed, but the response time is so slow that accurate laser power monitoring at each data sector is not possible.\nWhat is needed is a HAMR disk drive that can use a thermal sensor to accurately measure laser power during writing."}
{"text": "Running boards for vehicles are generally known, and are a popular accessory to attach to a vehicle because running boards both have functionality and can enhance the aesthetic appearance of a vehicle. Many types of running boards are designed specifically for trucks and sport-utility vehicles because these types of vehicles have a higher ride height compared to a car. However, the length of the wheel base for trucks and sport utility vehicles may vary. Trucks are manufactured with two conventional doors (commonly referred to as a “regular cab” design), two conventional doors and two second row half-doors (commonly referred to as an “extended cab” design), and four conventional doors (commonly referred to as a “quad cab” or “crew cab” design).\nVehicles having wheel bases of different lengths require running boards of different lengths to allow passengers entering and exiting the vehicle through the various doors to use the running boards. However, most conventional running boards having a tubular or flat shape are manufactured to be suitable for one particular wheel base.\nAccordingly, there exists a need for a running board which is suitable for use with different vehicles having wheel bases of different lengths."}
{"text": "Needleless access port valves are widely used in the medical industry for accessing an IV line and/or the internals of a patient or subject. Exemplary patents or publications that describe needleless access port valves are disclosed in Pub. Nos. 2002/0133124 A1 to Leinsing et al., Publication No. U.S. 2007/0191786 A1 and U.S. Pat. No. 6,871,838, both to Raines et al. The contents of each of the foregoing references are expressly incorporated herein by reference as if set forth in full. Generally speaking, prior art valves utilize a valve housing in combination with a moveable internal plug or piston to control the flow of fluid through a valve. The plug or piston may be moved by a syringe or a medical implement to open the inlet of the valve for accessing the interior cavity of the valve. When a fluid is delivered through the valve, fluid flow typically flows around the outside of the plug or piston in the direction towards the outlet. Upon removal of the syringe or medical implement, the plug or piston returns to its original position, either un-aided or aided by a biasing means, such as a spring or a diaphragm.\nIn some prior art valves, when the syringe or medical implement pushes the plug or piston, the plug or piston is pierced by an internal piercing device, such as a spike. The spike typically incorporates one or more fluid channels for fluid flow flowing through the pierced piston and then through the fluid channels in the spike. In yet other prior art valves, a self-flushing or positive flush feature is incorporated to push residual fluids confined inside the interior cavity of the valve to flow out the outlet when the syringe or medical implement is removed.\nConcerns with prior art valves include microbial colonization and sterilization. The former has to do with bacteria growth within the valve housing and the latter with the ability to clean the inlet of bacteria growth, such as wiping the inlet and/or the piston with a swab. Pub No. 2002/0133124 A1 to Leinsing et al., previously incorporated by reference, teaches the use of an antimicrobial agents, such as silver, silver oxide, or silver sulfadiazine, that “maybe included in the material forming the flex-tube piston or may be added to the outer surface of the piston as a coating” to minimize or eliminate bacteria growth. However, no information was provided on the effectiveness of the antimicrobial agents on bacterial growth when added as an admixture and cured within a piston. Nor was there information provided on how the antimicrobial agents kill the bacterium when the agents are cured or confined within the piston walls. Presumably, when cured within the piston, at most, only metals trapped along the outer surfaces of the piston are active agents."}
{"text": "1. Field of the Invention\nThe present invention relates to a reproducing control device and a reproducing control program that transmit a contents reproducing instruction to a contents reproducing device. Further, the present invention relates to the contents reproducing device that reproduces contents.\n2. Description of the Related Art\nPortable players (contents reproducing devices), which have HDD (hard disc drive) or flash memories for storing a plurality of contents (for example, music (tune) data) therein, are used. In the case where users reproduce contents stored in portable players and listen to music in their rooms, the portable players are connected to battery chargers (reproducing control devices) connected to amplifiers. The battery chargers have a function for supplying a power supply voltage to batteries of the portable players, and a function for transmitting audio signals from the portable players to the amplifiers. The audio signals reproduced in the portable players are transmitted via the battery chargers to the amplifiers, and the music is reproduced in the amplifiers. The portable players are normally provided with a resume function. The resume function is such that when a certain tune A which is being reproduced is stopped based on a users operation or when the power supply of the portable player is turned off, the tune A and its stop position (these are called resume information) are stored in a table in a RAM of the portable player, and when the user inputs a reproducing instruction into the portable player at the next time, the tune A whose reproduction is stopped at the last time is started to be played from the stop position.\nThe above type of portable player has the following problems. When the reproduction of the tune A is stopped with the portable player being connected to the battery charger, the portable player stores the resume information about the tune A in the table. When the portable player is removed from the battery charger and another tune B is played in the disconnected state, the resume information about the tune A stored in the table is deleted. When the reproduction of the tune B is stopped in the disconnected state, the resume information about the tune B is stored in the table. When the portable player is connected to the battery charger again and the reproduction is started, the resume information about the stopped tune B is stored in the table in the disconnected state. For this reason, the portable player starts the reproduction of the tune B. Even in the case where the user reproduces another tune B with the battery charger being disconnected from the portable player, when they are connected again, the tune A which was reproduced in the connected state at the last time is demanded to be automatically started to be reproduced. This demand is created due to such a circumstance that a tune to which the user listens in the disconnected state (namely, outside of the room) is different from a tune to which the user listens in the connected state (namely, in the room) depending on user's mood. In a conventional manner, however, since the resume information about the tune A is deleted at the time when the tune B is reproduced in the disconnected state, this demand cannot be realized."}
{"text": "New construction light pans that securely mount light fixtures and/or light trim modules are generally secured to ceiling joists or rafters located behind and supporting interior surface panels, also known in the industry as sheetrock, gypsum board and drywall. The pans are installed either before the ceiling panel is installed, or while access is still readily available from the back side of the panel. In contrast, when remodeling, a ceiling or wall panel is already in place and access from the back side of the panel is generally no longer available; therefore, all installation steps must be completed from the front side of the ceiling panel, generally through a cutout having a raw unfinished edge through which the remodel light fixture is installed. Additionally, a remodel light fixture usually must be secured to the panel cutout. The cutout in the panel, often lacks enough integrity and strength at the periphery of a cutout to support engagement of typical hardware such as formed spring steel clips used to secure the light fixture to the ceiling panel.\nPrior art remodel light fixtures include light fixtures having a can housing with a flange that projects radially from a lower end of the can and hardware to secure the can housing to the ceiling panel. For example, typical hardware includes jackscrew and flag combinations located peripherally around the can at a distance from the flange that allows the ceiling panel to be captured between the flange and flags. To install, the flags are swung inwardly tangential to the surface of the can, the distal end of the can is inserted into the hole through the ceiling panel, then the flange is held flush against the front side of the ceiling panel. Held in this position, the flags are located around the outer circumference of the can at a depth from the flange so that they are just above the back side of the ceiling panel. While the flange is held flush against the ceiling panel, each jackscrew is actuated and the associated flag swings radially outward and the flag is drawn toward the back side of the ceiling panel, capturing and compressing the panel between the flag and the flange. The light trim is then secured to the can housing. In some cases, the light receptacle coupled to the housing must be secured to the trim before the trim is secured to the housing.\nSuch a prior art remodel light fixture can be cumbersome to hold in place while securing and also requires separate installation of the light housing and trim module. Additionally, the raw, unfinished circumference of the cutout from a front side surface of a ceiling panel is often uneven, imperfect, and not structurally sound, requiring care that the diameter is proper and any crumbled areas or other imperfections do not prevent a flag from catching a solid portion of the back side of the panel or from crumbling a portion of the periphery of the cutout as the flag is swung out and drawn toward the flange.\nAnother prior art remodel light fixture includes a housing can and spring clips that can be manually actuated from the interior of the can and that extend radially outward from the outer diameter of the can. To install, the can is positioned within the hole in the ceiling panel so that the bottom edge of the can is flush with the front side of the ceiling panel. While the can is held in this position, the spring clips are manually actuated from the interior of the can so that the clips press into the panel thickness between the front and back side surfaces, securing the can in place. Next the light trim is secured to the housing can. In some fixtures, the light receptacle coupled to the housing must be secured to a trim module before the trim module is secured to the housing.\nInstallation of such remodel light fixtures can be cumbersome, requiring portions of the light assembly be held in place while securing and performing additional assembly of the housing, light receptacle, and trim. Additionally, the raw, unfinished circumference of the cutout from a front side surface of a ceiling panel such as drywall may crumble, give way, or be imperfect in opening size so as to prevent secure engagement by the spring clips that press into it.\nIt is therefore desirable to provide a remodel light assembly that installs easily, quickly, and securely in panels such as drywall, even when imperfections in the cutouts are present."}
{"text": "In recent years, sales for essential oils have exploded. Essential oils are usually oils which are derived from, or include certain essential components or essences of different plant substances. Such oils are generally ingested, topically applied, or are breathed in through various methods of diffusion or atomization.\nEssential oils, known as nature's living energy, are the natural, aromatic volatile liquids found in shrubs, flowers, trees, resins, fruit peels, rhizomes, roots, bushes, and seeds. The distinctive components in essential oils defend plants against insects, environmental conditions, and disease. They are also vital for a plant to grow, live, evolve, and adapt to its surroundings. Essential oils are extracted from aromatic plant sources via steam distillation, cold pressing, and other types of extraction and/or distillation. Essential oils are highly concentrated and far more potent than dry herbs. Other topically applied oils may include olive oil, almond oil, coconut oil, fatty acid oils, etc., and oils high in esters, such as jojoba oil, and waxes such as beeswax.\nWhile essential oils often have a pleasant aroma, their chemical makeup is complex and their benefits vast—which makes them much more than something that simply smells good. Historically, essential oils have played a prominent role in everyday life. With more than 200 references to aromatics, incense, and ointments throughout the Bible, essential oils are said to be used for anointing and healing the sick. Today, essential oils are used for aromatherapy, massage therapy, emotional health, personal care, nutritional supplements, household solutions, and much more.\nDiffusers for essential oils have been used to disperse the essential oils for breathing or to create a pleasant fragrance in a room or area. However, available diffusers for use with most essential oils are almost always unreliable with short service lives and high failure rate, and have to be refilled often. Additionally, most diffusers only operate for limited durations of a few hours before depleting the aromatic compounds. Many problems can be mitigated with meticulous care in maintaining the diffuser, which is beyond the capacity and patience of the average user.\nMany types of available diffusers use a piezo-electric ultrasonic transducer to agitate the water's surface, where essential oils reside, into water and essential vapor. The water has essential oils mixed in, some of which is transported along with the water into the air with a blower or fan. One of the greatest complaints with these types of diffusers is the low intensity of the essential oil aromas produced. The piezo transducer is located within the water and oil reservoir. When the water runs out, or is left to sit and evaporate, residual oils tend to collect on the piezo transducer. Accumulations of the oils on the piezo transducer reduce the effectiveness, and eventually cause failure of the diffuser, often in very short order. Users are instructed to carefully clean the diffuser after each use, including cleaning the transducer with a detergent. However, users often neglect this task or clean too vigorously, damaging or destroying the sensitive piezo transducer.\nOther diffusers use small fans or microblowers using traditional rotating fans. In small or micro fans, such microblowers are fairly inefficient, only supplying between about 5.5 and 60 Pa of pressure with a power consumption of between about 0.1 and 1.1 Watts with dimensions less than 25×25 mm. Any larger fans are not really microfans or microblowers. Such microfans also tend to have fairly short lives making diffusers run with such fans undesirable and inefficient. Additionally, such microfans are also unable to generate enough pressure through a small tube in order to push air through oil to create air saturated with oil for dispersion, but instead only generates sufficient air pressure to move air across the surface of oil, such as is shown in US Patent Pub. 2007/0138326."}
{"text": "A conventional sound absorbing material for an automobile comprising a felt material and a flame-resistant resin layer is known (for example, refer to Japanese Patent No. 3568936, (pages 3-5, FIG. 2), (Patent document 1)).\nThis sound absorbing material is obtained by applying a highly-viscous latex made of a flame-resistant resin as a coating to a felt material 2 so that a flame-resistant resin layer is formed on the surface of the felt material 2 and an independent fabric layer 8 remains whereby high sound absorbing properties are realized at medium and high sound ranges.\nAnother material including non-woven fabrics and a water-resistant film is known as a sound absorbing material as vehicle exterior equipment (for example, refer to Japanese Patent No. 3675359, (pages 2-4, FIG. 1) (Patent document 2)).\nThe sound absorbing material 11 as vehicle exterior equipment described in Patent Document 2 applies press molding in a state in which the water-resistant film 22 is closely attached to the surface of a fabric web 21 having a sheet shape and the two are mutually adhered.\nA material as vehicle exterior equipment including a non-woven fabric with one surface of a predetermined surface roughness and friction coefficient is known in other conventional technologies (for example, refer to JP2004-359066A, (page 9-10, FIG. 3) (Patent document 3)).\nAccording to Patent document 3, a first fiber aggregate and a second fiber aggregate are superimposed, heated and press molded. Thereby a material as vehicle exterior equipment including a layer 23 of the tire house side and a layer 22 of the tire side is manufactured.\nBy the way, it is desirable that foreign substances such as water, dust and dirt or the like do not adhere to a sound absorbing material used as vehicle exterior equipment. In particular, when a fender liner disposed in a tire house is used in a snowy area, once water, dirt, snow and ice or the like are adhered to the fender liner, it is not preferable that snow be further attached to the adhered substances or ice grows in the periphery of the adhered substances.\nIn the case of manufacturing a sound absorbing material described in Patent document 1, a highly-viscous fire-retarding material (latex) is coated on the felt member 2 to be heated and press molded. However, such a fire-retarding resin layer 5 hardly transmits heat when heated. Therefore, there is a problem in that molding by hot pressing is difficult.\nA sound absorbing material described in Patent document 2 applies press molding in a state in which the water-resistant film 22 is closely attached to the surface of the fiber web 21 having a sheet shape so that the water-resistant film 22 is adhered to the surface of the fiber web 21. The sound absorbing material is hereby press molded into a shape fitting in a vehicle body, but because the water-resistant film 22 hardly transmits heat, there is a problem in that moldability is poor. Furthermore, ready-made water-resistant films are expensive so that the manufacturing cost is increased.\nIn a material as vehicle exterior equipment described in Patent document 3, a first fabric aggregate and a second fabric aggregate are superimposed, heated and press molded. The average deviation of surface roughness and friction coefficient of one surface are limited to below the predetermined values. However, in this material as vehicle exterior equipment, the outward facing portion is comprised from a fiber aggregate so that there are limitations to reducing the fluff of the fabric."}
{"text": "Conventionally, there is an image pickup device such as a CMOS (Complementary Metal Oxide Semiconductor) image sensor in which a plurality of A/D converters (ADCs (Analog Digital Converters)) are connected to each output line to which pixels are connected for each column and high-speed data reading is achieved by using the plurality of ADCs (for example, see Patent Document 1).\nFurther, there is an image pickup device in which a plurality of comparators and counters are mounted on each output line, a voltage of a D/A converter (DAC (Digital Analog Converter)) is shifted by an arbitrary value, and a signal of high bit accuracy is read at high speed (for example, see Patent Document 2).\nFurther, there is an image pickup device in which noise is reduced and dynamic range is improved by performing A/D conversion twice on a read-out signal (for example, see Patent Document 3 and Patent Document 4)."}
{"text": "The present invention relates to a tilt-adjustable seat cushion for a vehicle seat having an underframe for receiving an upholstered seat part, this seat cushion having two frame parts which are arranged laterally of the underframe. The seat cushion includes a swiveling shaft for the underframe disposed in the frame parts and extends in transverse direction of the vehicle. Each frame part carries at least one row of teeth for interaction with two pinion gear teeth. An adjusting shaft carries the pinion gear teeth and has on one end a handwheel.\nIn a known seat cushion of the above type, shown in German Published Unexamined patent application No. 3,316,618, the adjusting shaft with the handwheel and the pinion gear teeth is disposed in the underframe and, for this purpose, on both sides, is guided in a respective longitudinal slot in the underframe. The rows of teeth that ascend in the direction of the swiveling shaft of the underframe, penetrate through openings in the floor of the underframe. For assembly, first the underframe is inserted into the frame parts. Then, the adjusting shaft is guided through one longitudinal slot in the underframe, through two recesses in the frame parts, of which one is in each case directly opposite one row of teeth, and then through the other longitudinal slot in the underframe, until two pinion gear teeth engage in the assigned row of teeth.\nBecause of the tolerances of the underframe and of the components for the tilt adjustment, the above-described assembly is not too easy. However, the assembly becomes significantly more difficult because it can only be carried out at the upholstered seat cushion, in other words, with the seat upholstery part that is firmly connected with the underframe. An objective of the present invention is to significantly facilitate the mounting of the underframe and of the tilt-adjusting mechanism in a seat cushion having an underframe for receiving an upholstered seat part.\nThis and other objectives are achieved in the present invention by providing a tilt-adjustable seat cushion for a vehicle seat having an underframe for receiving an upholstered seat part, the seat cushion having two frame parts which are arranged laterally of the underframe, swiveling shaft means for the underframe disposed in the frame parts that extends in transverse direction of the vehicle, each frame part carrying at least one row of teeth for interaction with at least two pinion gear teeth, and an adjusting shaft carrying the pinion gear teeth having at one end a hand wheel, with bearing block means. The adjusting shaft is disposed in the bearing block means, this bearing block means being connected to and forming a separate assembly with the frame parts and the adjusting shaft. This separate assembly is then connected with the underframe by fastening means. In certain preferred embodiments, the fastening means are screws or clips.\nBy disposing the adjusting shaft in a separate bearing block that is separate from the underframe, the tilt-adjusting mechanism comprising the frame part and the adjusting shaft can be assembled completely and only then be connected with the upholstered seat cushion. Because the underframe and the seat upholstery part are absent during the assembly of the tilt-adjusting mechanism, there is more free space and better accessibility so that assembly times are considerably shortened. After the pre-assembly of the tilt-adjusting mechanism is completed, the upholstered underframe is placed without difficulty and connected with the adjusting mechanism.\nIn certain preferred embodiments, the bearing block or at lest the bearing points of the adjusting shaft in the housing block are made of plastic, to prevent the development of noise caused by rattling, etc.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "The method most commonly used to produce iron oxides used as pigments employs, as its source material, pickling liquids that contain ferrous sulfate or chloride originating from the iron industry or from the process for producing titanium dioxide.\nThe acid aqueous solution of ferrous salts is first neutralized by dissolving iron scrap and is then treated in conditions of oxidation with sodium hydroxide to precipitate FeOOH nuclei which are then pumped into a reactor that contains iron scrap and is maintained in air stream.\nThe resulting ferric sulfate (chloride) hydrolyzes forming FeOOH or Fe2O3; the sulfuric or hydrochloric acid that is released reacts with the iron to form ferrous sulfate or chloride, which are then oxidized to ferric salts. The reaction time varies from a few days to several weeks, depending on the reaction conditions and on the type of pigment that is sought.\nThe advantage of the method, with respect to others, is the limited use of alkali and of ferrous sulfate or chloride. The small quantity of ferrous salt that is required initially is renewed continuously during the process by the dissolving of the iron by the sulfuric or hydrochloric acid released in the reaction.\nThe disadvantage of the method is the difficulty in eliminating, even after thorough washing, the impurities of sulfate and chloride anions that are present in the oxides, which have a negative effect on the quality of the pigments.\nFor example, in order to reduce these anions to values that are acceptable for the production of high-quality red pigments, it is necessary to treat the precipitated oxides with concentrated solutions of NaOH (U.S. Pat. No. 5,614,012).\nGB 1226876 describes a method for producing highly pure FeOOH suitable for producing ferrites for use in electronic devices, wherein electrolytic iron with average dimension between 20 and 140 microns is reacted, in conditions of oxidation with air made to flow at high speed in order to maintain a uniform aqueous suspension of the iron particles, with an acid chosen among sulfuric acid, hydrochloric acid, nitric acid and acetic acid, used at a molar concentration of less than 0.01 and in a molar ratio with the iron of more than 0.02 and preferably between 0.26 and 0.55. Iron is used in an amount not exceeding 25 g/l and the weight ratio between the solution and the iron is at least 40.\nThe reaction temperature is between 50 and 70° C.: at temperatures above 70° C., there is an undesirable production of oxides such as spinel, which also form at temperatures below 70° C. if the iron concentration is higher than 25 g/l.\nAt temperatures below 50° C., the oxide particles that form are too fine and difficult to filter and wash in order to achieve values of impurities due to acid radicals of less than 0.1% by weight.\nThe sought dimensions of FeOOH are a few microns in length and more than 0.3 and 0.1 microns in width and thickness, respectively.\nIf the concentration of the acid is too high (more than 0.25 mol in the case of sulfuric acid), the FeOOH yield decreases even considerably due to the dissolving of iron ions in the mother liquor. The productivity of the method is 20-26 g of FeOOH per liter of suspension per hour."}
{"text": "JPEG as described in W. Pennebaker and J. Mitchell, “JPEG still image data compression standard,” Kluwer Academic Publishers, 1993, (hereinafter “reference [1]”), G. Wallace, “The JPEG still-image compression standard,” Commun. ACM, vol. 34, pp. 30-44, April 1991 (hereinafter “reference [2]”), is a popular DCT-based still image compression standard. It has spurred a wide-ranging usage of JPEG format such as on the World-Wide-Web and in digital cameras.\nThe popularity of the JPEG coding system has motivated the study of JPEG optimization schemes—see for example J. Huang and T. Meng, “Optimal quantizer step sizes for transform coders,” in Proc. IEEE Int. Conf. Acoustics, Speech and Signal Processing, pp. 2621-2624, April 1991 (hereinafter “reference [3]”), S. Wu and A. Gersho, “Rate-constrained picture-adaptive quantization for JPEG baseline coders,” in Proc. IEEE Int. Conf. Acoustics, Speech and Signal Processing, vol. 5, pp. 389-392, 1993 (hereinafter “reference [4]”), V. Ratnakar and M. Livny, “RD-OPT: An efficient algorithm for optimizing DCT quantization tables”, in Proc. Data Compression Conf., pp. 332-341, 1995 (hereinafter “reference [5]”) and V. Ratnakar and M. Livny, “An efficient algorithm for optimizing DCT quantization,” IEEE Trans. Image Processing, vol. 9 pp. 267-270, February 2000 (hereinafter “reference [6]”), K. Ramchandran and M. Vetterli, “Rate-distortion optimal fast thresholding with complete JPEG/MPEG decoder compatibility,” IEEE Trans Image Processing, vol. 3, pp. 700-704, September 1994 (hereinafter “reference [7]”), M. Crouse and K. Ramchandran, “Joint thresholding and quantizer selection for decoder-compatible baseline JPEG,” in Proc. IEEE Int. Conf. Acoustics, Speech and Signal Processing, pp. 2331-2334, 1995 (hereinafter “reference [8]”) and M. Crouse and K. Ramchandran, “Joint thresholding and quantizer selection for transform image coding: Entropy constrained analysis and applications to baseline JPEG,” IEEE Trans. Image Processing, vol. 6, pp. 285-297, February 1997 (hereinafter “reference [9]”). The schemes described in all of these references remain faithful to the JPEG syntax. Since such schemes only optimize the JPEG encoders without changing the standard JPEG decoders, they can not only further reduce the size of JPEG compressed images, but also have the advantage of being easily deployable. This unique feature makes them attractive in applications where the receiving terminals are not sophisticated to support new decoders, such as in wireless communications.\nQuantization Table Optimization\nJPEG's quantization step sizes largely determine the rate-distortion tradeoff in a JPEG compressed image. However, using the default quantization tables is suboptimal since these tables are image-independent. Therefore, the purpose of any quantization table optimization scheme is to obtain an efficient, image-adaptive quantization table for each image component. The problem of quantization table optimization can be formulated easily as follows. (Without loss of generality we only consider one image component in the following discussion.) Given an input image with a target bit rate Rbudget, one wants to find a set of quantization step sizes {Qk:k=0, . . . , 63} to minimize the overall distortion\n                    D        =                              ∑                          n              =              1                                      Num              ⁢              _              ⁢              Blk                                ⁢                                    ∑                              k                =                0                            63                        ⁢                                          D                                  n                  ,                  k                                            ⁡                              (                                  Q                  k                                )                                                                        (        1        )            subject to the bit rate constraint\n                    R        =                                            ∑                              n                =                1                                            Num                ⁢                _                ⁢                Blk                                      ⁢                                          R                n                            ⁡                              (                                                      Q                    0                                    ,                  …                  ⁢                                                                          ,                                      Q                    63                                                  )                                              ≤                      R            budget                                              (        2        )            where Num_Blk is the number of blocks, Dn,k(Qk) is the distortion of the kth DCT coefficient in the nth block if it is quantized with the step size Qk, and Rn(Q0, . . . , Q63) is the number of bits generated in coding the nth block with the quantization table {Q0, . . . , Q63}.\nSince JPEG uses zero run-length coding, which combines zero coefficient indices from different frequency bands into one symbol, the bit rate is not simply the sum of bits contributed by coding each individual coefficient index. Therefore, it is difficult to obtain an optimal solution to (1) and (2) with classical bit allocation techniques. Huang and Meng—see reference [3]—proposed a gradient descent technique to solve for a locally optimal solution to the quantization table design problem based on the assumption that the probability distributions of the DCT coefficients are Laplacian. A greedy, steepest-descent optimization scheme was proposed later which makes no assumptions on the probability distribution of the DCT coefficients—see reference [4]. Starting with an initial quantization table of large step sizes, corresponding to low bit rate and high distortion, their algorithm decreases the step size in one entry of the quantization table at a time until a target bit rate is reached. In each iteration, they try to update the quantization table in such a way that the ratio of decrease in distortion to increase in bit rate is maximized over all possible reduced step size values for one entry of the quantization table. Mathematically, their algorithm seeks the values of k and q that solve the following maximization problem\n                              max          k                ⁢                              max            q                    ⁢                                                                      -                  Δ                                ⁢                                                                  ⁢                D                            ⁢                              |                                                      Q                    k                                    ->                  q                                                                                    Δ                ⁢                                                                  ⁢                R                            ⁢                              |                                                      Q                    k                                    ->                  q                                                                                        (        3        )            where ΔD|Qk→q and ΔRQk→q are respectively the change in distortion and that in overall bit rate when the kth entry of the quantization table, Qk, is replaced by q. These increments can be calculated by\n                                              ⁢                                                            Δ                ⁢                                                                  ⁢                D                            ⁢                              |                                                      Q                    k                                    ->                  q                                                      =                                          ∑                                  n                  =                  1                                                  Num                  ⁢                  _                  ⁢                  Blk                                            ⁢                              [                                                                            D                                              n                        ,                        k                                                              ⁡                                          (                      q                      )                                                        -                                                            D                                              n                        ,                        k                                                              ⁡                                          (                                              Q                        k                                            )                                                                      ]                                              ⁢                                          ⁢                                          ⁢          and                                    (        4        )                                                      Δ            ⁢                                                  ⁢            R                    ⁢                      |                                          Q                k                            ->              q                                      =                              ∑                          n              =              1                                      Num              ⁢              _              ⁢              Blk                                ⁢                      [                                                            R                  n                                ⁡                                  (                                                            Q                      0                                        ,                    …                    ⁢                                                                                  ,                    q                    ,                    …                    ⁢                                                                                  ,                                          Q                      63                                                        )                                            -                                                R                  n                                ⁡                                  (                                                            Q                      0                                        ,                    …                    ⁢                                                                                  ,                                          Q                      k                                        ,                    …                    ⁢                                                                                  ,                                          Q                      63                                                        )                                                      ]                                              (        5        )            The iteration is repeated until |Rbudget−R(Q0, . . . , Q63)|≦ε where ε is the convergence criterion specified by the user.\nBoth algorithms aforementioned are very computationally expensive. Ratnakar and Livny—see references [5] and [6]—proposed a comparatively efficient algorithm to construct the quantization table based on the DCT coefficient distribution statistics without repeating the entire compression-decompression cycle. They employed a dynamic programming approach to optimizing quantization tables over a wide range of rates and distortions and achieved a similar performance as the scheme in reference [4].\nOptimal Thresholding\nIn JPEG, the same quantization table must be applied to every image block. This is also true even when an image-adaptive quantization table is used. Thus, JPEG quantization lacks local adaptivity, indicating the potential gain remains from exploiting discrepancies between a particular block's characteristics and the average block statistics. This is the motivation for the optimal fast thresholding algorithm of—see reference [7], which drops the less significant coefficient indices in the R-D sense. Mathematically, it minimizes the distortion, for a fixed quantizer, between the original image X and the thresholded image {tilde over (X)} given the quantized image {circumflex over (X)} subject to a bit budget constraint, i.e.,\n                              min          ⁢                      ⌊                                          D                ⁡                                  (                                      X                    ,                                          X                      ~                                                        )                                            |                              X                ^                                      ⌋                    ⁢                                          ⁢          subject          ⁢                                          ⁢          to          ⁢                                          ⁢                      R            ⁡                          (                              X                ~                            )                                      ≤                  R          budget                                    (        6        )            \nAn equivalent unconstrained problem is to minimize\n                              J          ⁡                      (            λ            )                          =                              D            ⁡                          (                              X                ,                                  X                  ~                                            )                                +                      λ            ⁢                                                  ⁢                          R              ⁡                              (                                  X                  ~                                )                                                                        (        7        )            \nA dynamic programming algorithm is employed to solve the above optimization problem (7) recursively. It calculates J*k for each 0≦k≦63, and then finds k* that minimizes this J*k, i.e., finding the best nonzero coefficient to end the scan within each block independently. The reader is referred to reference [7] for details. Since only the less significant coefficient indices can be changed, the optimal fast thresholding algorithm—see reference [7]—does not address the full optimization of the coefficient indices with JPEG decoder compatibility.\nJoint Thresholding and Quantizer Selection\nSince an adaptive quantizer selection scheme exploits image-wide statistics, while the thresholding algorithm exploits block-level statistics, their operations are nearly “orthogonal”. This indicates that it is beneficial to bind them together. The Huffman table is another free parameter left to a JPEG encoder. Therefore, Crouse and Ramchandran—see references [8] and [9]—proposed a joint optimization scheme over these three parameters, i.e.,\n                                          min                          T              ,              Q              ,              H                                ⁢                                    D              ⁡                              (                                  T                  ,                  Q                                )                                      ⁢                                                  ⁢            subject            ⁢                                                  ⁢            to            ⁢                                                  ⁢                          R              ⁡                              (                                  T                  ,                  Q                  ,                  H                                )                                                    ≤                  R          budget                                    (        8        )            where Q is the quantization table, H is the Huffman table incorporated, and T is a set of binary thresholding tags that signal whether to threshold a coefficient index. The constrained minimization problem of (8) is converted into an unconstrained problem by the Lagrange multiplier as\n                              min                      T            ,            Q            ,            H                          ⁢                  [                                    J              ⁡                              (                λ                )                                      =                                          D                ⁡                                  (                                      T                    ,                    Q                                    )                                            +                              λ                ⁢                                                                  ⁢                                  R                  ⁡                                      (                                          T                      ,                      Q                      ,                      H                                        )                                                                                ]                                    (        9        )            \nThen, they proposed an algorithm that iteratively chooses each of Q, T, H to minimize the Lagrangian cost (9) given that the other parameters are fixed.\nJPEG Limitations\nThe foregoing discussion has focused on optimization within the confines of JPEG syntax. However, given the JPEG syntax, the R-D performance a JPEG optimization method can improve is limited. Part of the limitation comes from the poor context modeling used by a JPEG coder, which fails to take full advantage of the pixel correlation existing in both space and frequency domains. Consequently, context-based arithmetic coding is proposed in the literature to replace the Huffman coding used in JPEG for better R-D performance."}
{"text": "1. Field of the Invention\nThe present invention generally relates to debugger systems for computer software.\n2. Description of the Prior Art\nDuring the coding of computer programs, the inadvertent production of errors, or bugs, is typically unavoidable. Debugging systems help identify, and thus eliminate, these bugs from a computer program.\nOne tool used with debuggers are breakpoints. Breakpoints are stopping locations in a computer program under test. The computer program will run until it reaches a predetermined breakpoint address; control is then returned to the debugger. At this point, the debugger can examine the memory locations used by the computer program under test. By examining these memory locations, the user can locate and eliminate bugs in the computer program under test.\nThere are difficulties with breakpoints, however. For example, some computer systems are parts of systems in which the computer system under test provides and receives data from another part of the system within a certain time restriction. The breakpoints can thus cause timing problems. Sometimes it is not feasible for the computer user to be around during the testing of the computer system. For example, if a breakpoint is set at a location where it is unlikely for the program to hit, it would not make sense for the computer user to sit at the computer screen waiting for such an unlikely event. Additionally, some testing of the computer system can occur late at night when the software engineers are not available.\nOne way to avoid these problems is through the use of tracepoints. Tracepoints are similar to breakpoints, but unlike breakpoints, they do not stop the program indefinitely. Once the program reaches a tracepoint, data in predetermined memory locations are stored. Control is then returned to the computer program under test. After the computer program has run, a xe2x80x9chex dumpxe2x80x9d of all the data obtained by the tracepoints is provided to the user. The user sifts through the hexadecimal representation of the data in order to determine the bugs in the computer system. It is desired to have an improved tracepointing system for examining a computer program under test.\nOne embodiment of the present invention relates to the use of programming abstractions to represent the tracepoint output as well as the data to collect at each tracepoint event. This makes it easier for the user of the debugging system to collect and evaluate the tracepoint event data. The user does not have to sift through the hexadecimal representations in order to determine the locations of and contents of memory. For example, if a variable, X, is used in the source code of the computer program under test, the debugger system can interpret the user\"\"s command xe2x80x9ccollect Xxe2x80x9d at a given tracepoint to collect the data stored in the memory location associated with the variable, X. Later, an instruction to xe2x80x9cprint Xxe2x80x9d will cause the debugging software to display the value of X. The user need not look up the address of variable X, then look at the contents of that address; the debugger system does this automatically for the user.\nIn a preferred embodiment, the system can collect data from the stack and then produce an output indicating what subroutine called which, and with what arguments, at the time of the tracepoint event.\nIn a preferred embodiment, the debugger system uses a programming abstraction, such as a variable used in the source code of the program under test, and checks the address and symbol table created during the compiling of the computer program under test to determine some indication of the addresses of the required data. This indication is sent to a debugger agent (a subprogram running on the target system that provides services for the debugger on the host computer) along with the computer program under test. The debugger agent in the target system collects the desired data and stores the data into a buffer when a tracepoint is executed. After the computer program under test has run, the debugger system can send a request to the debugger agent on the target to obtain the data in the buffer. The debugger can then represent the data using programming abstractions. In a preferred embodiment, both the collection and display representations of the data uses the symbol and address table created by the compiler for the compiled program under test.\nIn another embodiment of the present invention, the system includes bytecode expressions that are run as part of the tracepoint collecting of data at the target system. The target system can evaluate expressions in order to determine the data to collect. For example, consider an array Q[n] which has 100 elements. If the user wants to collect the data in the location Q[X+Y], where X and Y are both program variables, the present invention can evaluate X+Y in order to determine the address location of the array value Q[X+Y] in order to obtain this value and store it in the buffer. In this way, less data need be stored. For example, if only the element Q[X+Y] is required, the system need not store all 100 elements of the array each time the tracepoint is hit. This is an important consideration because the data storage capacity of the buffer at the target system is an important limiting factor. A given tracepoint could be hit a large number of times. By using bytecode expressions, the buffer is not filled with unwanted data. and thus filled prematurely.\nHaving the arithmetic expression evaluated at the target system rather than at the debugger is important. In breakpoint systems, typically the evaluation of expressions is done at the debugger. This is undesirable when using tracepoints that operate in real time because the arithmetic calculations of the expressions at the target system can occur much faster than the data transfer between the target system and the host system. For the same reason, the debugger agent stores the data into a buffer, rather than sending the data directly to the host system at the time of a tracepoint event."}
{"text": "Many disabled and elderly persons are inhibited or precluded from playing golf due to various aspects of their conditions. Some may have trouble getting into position to swing a golf club, such as those who rely on a wheelchair for mobility. Others may not be able to effectively grip and/or move their body to swing the club, such as persons suffering from certain paralysis. Still others may be physically able, but lack the cognitive ability to swing a golf club in a traditional manner, such as some mentally disabled or autistic persons.\nAdditionally, many persons seek to improve their golf swing through the use of training apparatuses. Training apparatuses physically manipulate a person's movement or the movement of their club to teach certain swing mechanics. Many persons learn more effectively by witnessing visual demonstrations of certain techniques including, but not limited to, the pendulum-like swing motion often used in chipping and putting."}
{"text": "Force control in robotics is conventionally implemented using either impedance control or admittance control. One example of an impedance control feedback loop is illustrated in FIG. 1. With impedance control, positions of the joints of the robot are inputted into the controller and joint torques for controlling movement of the robot are outputted and applied. In other words, the impedance controller determines position and applies (or commands) force/torque. In FIG. 1, the impedance controller applies specific joint torques to the joints. If the robot experiences external force acting on one of the joints, for example, the impedance control system does not calculate or measure such force. Instead, the impedance controller merely re-determines the robot position and re-calculates the requisite force to be applied.\nConventional impedance control may provide stable control when contacting rigid environments and may provide a light feel when engaging soft environments. However, impedance control can give the robot an unstable loose feel and may introduce errors when interacting with stiff virtual constraints, such as haptic boundaries, which limit movement of the robot.\nAdmittance control, on the other hand, is the inverse of impedance control. One example of an admittance control feedback loop is illustrated in FIG. 2. With admittance control, rather than determining position and commanding force, the controller instead determines applied force/torque and commands position. A force-torque sensor or joint torque measurements are used to detect input force to the system. Based on the detected input force, and knowing a current position of the joints based on measured joint angles, the admittance controller commands a new position of the joints by applying determined joint torques to move the joints accordingly.\nConventional admittance control can give the robot stable rigid feel and may reduce errors when interacting with virtual constraints, such as haptic boundaries. However, a robot subject to admittance control may feel heavy to a user and may overreact when contacting rigid environments. As significantly, using a single admittance controller that utilizes either the force/torque sensor or the joint torques to measure external force(s) acting on one or more of the joints provides significant challenges. Mainly, when the robot experiences such external forces, the location(s) (e.g., the joint(s)) to which the external forces are applied are unknown thereby potentially resulting in undesired dynamic behavior of the robot."}
{"text": "1. Field of the Invention\nThis invention relates to a circuit arrangement for producing high voltages by means of a voltage multiplier, e.g. of the Cockcroft-Walton type, that is fed from an alternating current source.\n2. Description of the Related Technology\nIt is known per se to provide high-voltages e.g. for discharge tubes, including X-ray tubes, by means of voltage multipliers of the kind concerned.\nIn the case of X-ray tubes it is generally desired to earth the anode because this makes it easier to obtain the necessary dispersion of heat from the anode. In order to avoid extremely high potentials in relation to earth, it is a frequently chosen solution, however, to supply the anode and cathode of the X-ray tube from high-voltage sources of opposite polarities each at half of the full anode-cathode voltage of the tube.\nIn either case, the filament of the tube will attain a high voltage relative to earth and due to the fact that the magnitude of the filament current determining the intensity of the X-radiation shall, moreover, generally be adjustable, certain difficulties arise, in particular as regards fulfilment of the demands for insulation, whereby the means for supplying and controlling the filament current become excessively voluminous and heavy, meaning a considerable increase of the total volume and weight of the X-ray tube. This is particularly inadvantageous in case of mobile X-ray devices for radiological tests of materials.\nThis applies also to other parts of the tube, e.g. a bias grid, the voltage of which deviates by only a small amount from the cathode voltage of the tube.\nSimilar problems occur also in other devices than X-ray tubes, which are supplied from a circuit arrangement of the type in question, and where it is desirable to produce small amplitude voltages lying at the voltage level of the high voltage."}
{"text": "The invention relates to controllably dissolving a composite.\nControlled release of medication in vivo is the subject of much research. Various methods and release agents have been suggested, tested and marketed. Calcium sulfate has been utilized as filler for bone cavities as it is capable of being spontaneously adsorbed and replaced by bone. Calcium sulfate, formed from the hemihydrate, has been used as a controlled release agent alone for the filling of bone cavities and in combination with additives such as medicaments and pesticides. As a carrier for medicaments, it has been useful in vivo because it is biocompatible and is progressively resorbed by the body, thereby eliminating the need for secondary surgical procedure.\nOne application for a calcium sulfate controlled release agent is the local delivery of medicaments in vivo. The ideal characteristics of a local medicament delivery system are (1) biodegradability, (2) biocompatibility, (3) prolonged pharmaceutical release (e.g., over a period of at least 4 to 6 weeks), (4) reproducibility, (5) predictable pharmacokinetics, and (6) controllability.\nOne of the disadvantages to the use of calcium sulfate as a carrier is that, for some medicaments, the medicament is eluted from the calcium sulfate matrix at too rapid of a rate."}
{"text": "(1) Field of the Invention\nThe present invention relates generally to aligning wafers, and in particular, to aligning wafers independent of the layers formed therein in the manufacture of integrated circuits.\n(2) Description of the Related Art\nAlignment marks and aligning wafers with respect to those marks are an important part of the process of manufacturing semiconductor devices and integrated circuits. As is known in the art, integrated circuits are fabricated by patterning a sequence of masking layers, and the features on successive layers bear a spatial relationship to one another. Thus, as a part of the fabrication process each level must be aligned to the previous levels. Alignment of one pattern layer to previous layers is done with the assistance of special alignment patterns designed on to each mask level. When these special patterns are aligned, it is assumed that the remainder of the circuit patterns are also correctly aligned. Since each layer must have alignment marks for proper registration with respect to the next layer, each alignment or registration then becomes layer dependent.\nThe tools that are used to pattern the various layers in a wafer are known as photomasks or masks and reticles. The patterns on the mask or the reticle are defined by a combination of opaque and translucent areas. A light source through the mask or the reticle projects the patterns onto the surface of a wafer, and depending upon the material that is being exposed to the light, the pattern is transferred on to the surface where the light arrives or not. A mask contains patterns that can be transferred to an entire wafer in one exposure. A reticle, on the other hand, contains a pattern image which must be stepped and repeated in order to expose an entire substrate.\nThe adjustment of the image of the mask being exposed to the previously produced patterns was originally performed by human operators, who compared the image locations under a microscope and adjusted the position of the mask to bring it into alignment with wafer patterns. Decreasing feature sizes, and the increasing number of alignments per wafer with step-and-repeat projection aligners, have been the impetus for developing automatic alignment systems. The principle of one type of automatic alignment procedure is illustrated on page 476 of S. Wolf's book, \"Silicon Processing for the VLSI Era,\" vol. 2, Lattice Press, Sunset Beach, Calif., 1990. Alignment marks consisting of two rectangular patterns, each set at a 45.degree. angle to the directions of the motion an xy-stage table are fabricated on a wafer. Two corresponding rectangular patterns are located on the reticle of an optical aligner, and their image is projected onto the wafer. The superimposed alignment target and the reticle image are reflected back into the main optical element of the aligner, and then into an on-axis microscope. The image from the microscope is focused onto the face of a TV camera, and is subsequently digitized into a form that can be analyzed by a computer. When alignment is achieved, a signal is generated. The 45.degree. orientation of the alignment marks makes it possible to obtain both x an y registration information from the horizontal scan-lines of the video camera. The relative position of the wafer marks with respect to the reticle windows determines the registration of the two images.\nRegistration and alignment in wafer steppers are performed globally and locally. Global alignment performs rotational and translational alignment of the entire wafer, while local alignment provides alignment to a target within a particular region of a pattern image on the wafer. Global alignment is usually accomplished at a remote alignment station before a wafer is moved under the projection lens for exposure. An apparatus for global aligning is disclosed in U.S. Pat. No. 4,046,985, though for a different purpose other than for aligning masks. The aforementioned apparatus aligns the wafer to a fixed reference position, inverts the aligned wafer to expose its backside and transfers it with controlled motion to a set position under a laser apparatus. The laser beam scans the backside of the wafer to create in a kerf area between chip sites on the wafer an easily breakable cut for the purpose of subsequently separating and removing the chips. This method of aligning differs from the method that will be disclosed in the present invention.\nU.S. Pat. No. 3,752,589 discloses a method of aligning the pattern of a photomask on one side of a wafer to the patterns placed on the underside of the wafer. The alignment of the mask with respect to the wafer is achieved by optically superimposing the images present on the mask and on the underside of the wafer and adjusting the mask relative to the wafer until the relative positions of the combined images are corrected to a predefined set of conditions. This method requires two viewing apparatus: one for the mask facing the wafer and the other on the opposing side facing the underside of the wafer, which is superfluous in the present invention as disclosed later in a preferred embodiment.\nCade, in U.S. Pat. No. 4,534,804, on the other hand, teaches a method of forming a laser alignment mark such that same mark now extending, after heat treatment, to both the front and the back of the wafer can be used to align photoresist photomasks to both the front and the back side of a silicon wafer. The alignment mark that is formed is actually a defect planted into the wafer by means of a laser beam. The wavelength of the laser beam is chosen such that it passes through the lightly doped wafer without absorption but is absorbed by a following heavily doped semiconductor layer to generate therein heat and resulting defect. The semiconductor wafer is then heated to cause the defects to migrate through a lightly doped epitaxial layer to the front surface thereof in which there is formed an identically positioned image of the mark scribed on the back side. It will be appreciated, however, that this method is limited to contiguous semiconductor materials without any intervening other types of layers.\nGenerally, the process of aligning masks in fabricating semiconductor devices in a substrate, and, subsequently, in \"metallizing\" or wiring the devices together to form integrated circuits on the substrate require different considerations. Up to the level where devices are fabricated, the alignment of masks is accomplished by projecting infrared rays from the underside of the wafer while observing the patterns by an infrared microscope. But this method is not applicable when metallized layers are present because metallic films are opaque to infrared rays. Similarly, the method of U.S. Pat. No. 4,534,804 cannot applied since the metallized layer would preclude the migration of a laser defect through it.\nThe process of metallization, or \"personalization\" requires the patterning of the metal layers to form the desired circuit pattern. This is accomplished by masking the metal layer with a photosensitive emulsion and then positioning a photographic mask thereon. The emulsion is next exposed to ultraviolet rays though the photographic mask and then the emulsion is developed and the unexposed emulsions is washed away with solvent. The exposed metal areas are then etched to form wiring patterns corresponding to those on the mask. Certain predetermined areas on the metal layer do also contain alignment marks for registration with the next level of metal layer to be deposited.\nBefore the next metal layer is deposited, however, an interlevel layer of a dielectric insulator is first blanket deposited over the wiring layer. An interlevel layer usually will form a relatively rough topography conforming to the geometrical features of the underlying layer. Since the depth-of-field limitations of submicron optical-lithography tools require surfaces to be planar within .+-.0.5 micrometers (See S. Wolf and R. N. Tauber, \"Silicon Processing for the VLSI Era,\" vol. 2, Lattice Press, Sunset Beach, Calif., 1990, p. 203), planarization, as is well known in the field, of the dielectric layer will also be required if optical lithography is to be usable for fabricating integrated circuits with submicron feature sizes of today. Conventionally, planarization is performed in any number of ways including mechanical polishing, or a combination of mechanical and chemical polishing, called CMP. It will be appreciated that as more layers of metal and interlevel dielectric are deposited, the surface topography of each layer will vary because of the cumulative effect of the number of underlying features that are disposed on top of each other. Hence, the photographic alignment of masks of each layer to the preceding layer will depend upon the nature of the planarized surfaces of each layer, and therefore will vary accordingly. A method is disclosed in this invention where the layer dependency of alignment is eliminated.\nTo appreciate the layer dependency of alignment and its effects thereof, FIGS. 1a and 1b show a conventional system where wafer (20) is disposed between a movable xy-stage (30) and an alignment source (10). The rays emanating from source (10) are shown by numeral (15) in both FIGS. 1a and 1b. In FIG 1a, wafer (20) is provided with one level of metal (40), or wiring layer, and in FIG. 2a, with two levels of metal, where the second wiring layer is denoted by numeral (50). First and second metal layers, (40) and (50), respectively, are separated by an interlevel dielectric layer (45), and second layer (50) has on it dielectric layer (55). As is commonly used in the manufacture of semiconductor wafers, the interlevel dielectric layers (45) and (55) shown in FIGS. 1a and 1b are spin-on-glass, or SOG. For the reasons given earlier, both of the dielectric layers are subjected to planarization by means of chemical-mechanical polishing before the respective metals layers are deposited thereon. However, it will be noted that the roughness or non-planarity of second dielectric layer (55) is more pronounced than that of the first layer (45). This is because, as the number of metal interconnect levels are increased, the stacking of additional layers on top of one another produces a more and more rugged topography as mentioned earlier. Although polishing creates more planarity, nonetheless, the remaining non-planarity varies from one layer to another depending upon the cumulative effect on the topography of the underlying layer. As a result, the patterns on layer (50) visible through the second interlevel layer (55), for example, are more diffused and not as clear as the patterns on layer (40) are through layer (45). Furthermore, the depth-of-field varies from layer to layer on the optical-lithography tools that are used for aligning masks over different layers. Consequently, the signals generated by the edges of the patterns on metal layers vary depending upon the planarity of the layer over which a mask is placed. In other words, the signals for aligning masks over patterns and alignment marks on a wafer are layer dependent, and are sometimes weak, and at other times not distinct so as to cause misalignment and therefore product defects. What is needed, therefore, is a method whereby the aligning and alignment signals are independent of a particular layer over which a mask is positioned and that the alignment signal strength is invariant throughout the manufacture of a semiconductor substrate."}
{"text": "1. Field of the Invention\nThe present invention relates to an actuator for lifting device and more particularly, to such an actuator, which comprises a linking rod assembly for moving a lifting device, and a transmission gear set and a brake module arranged at one side of the linking rod assembly and rotatable by a motor to move the linking rod assembly in extending out or retracting an extension tube.\n2. Description of the Related Art\nA lifting device is a mechanism for changing the elevation of an object. Different lifting devices differently configured to fit different application requirements.\nFIG. 11 illustrates an actuator A for use in a lifting device. As illustrated, the actuator A comprises an electric motor A1 having a worm A11 fixedly connected to the output shaft thereof, a worm gear A2 meshed with the worm All, a spiral spring A3 sleeved onto the gear shaft of the worm gear A2, and a retractable rod A4 coupled to the worm gear A2. According to this design, the retractable rod A4 and the worm gear A2 are on the same axis. When starting the electric motor A1, the worm A11 is driven to rotate the worm gear A2, the retractable rod A4 is moved axially subject to rotation of the worm gear A2, and the spiral spring A3 is wound tightly on the periphery of the gear shaft of the worm gear A2.\nThe aforesaid prior art actuator provides a braking effect to prevent an accident, however it still has drawbacks as follows:\n1. Because the retractable rod A4, the spiral spring A3 and the worm gear A2 are arranged on one same path, the worm bear A2 occupies a part of the reciprocating path of the retractable rod A4. Thus, the length of the retractable rod A4 is limited, lowering the performance of the actuator.\n2. The worm gear A2 is directly connected in series to one end of the retractable rod A4. When the retractable rod A4 receives an external pressure, the pressure will be directly transferred by the retractable rod A4 to the worm gear A2. If the worm gear A2 is held down by the spiral spring A3 at this time, direct transfer of the external pressure by the retractable rod A4 to the worm gear A2 may damage the worm gear A2 and the spiral spring A3."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for producing powder and a fluidized bed pulverizing apparatus.\n2. Description of the Related Art\nA toner used for an electrophotographic image forming apparatus is formed of fine particles having relatively uniform particle sizes of micron order. As an apparatus for producing such fine particles (powder) of micron order, a fluidized bed pulverizing apparatus (also called as an air flow pulverizing apparatus) is known. The fluidized bed pulverizing apparatus is constituted with a pulverization chamber (fluidized bed container), in which pulverization of a powder material is performed by allowing the powder material to collide against each other, a plurality of fluid jetting nozzles for jetting fluid in the pulverization chamber so as to entrain the powder material in the fluid, followed by colliding against each other so that the powder material entrained therein also collide each other, and then forming a fluidized bed in which the powder material further collide and are pulverized, and a centrifugal classification rotor which classifies the finely-pulverized powder, and is provided at the upper part of the pulverization chamber. In a typical fluidized bed pulverizing apparatus, the powder material supplied into the pulverization chamber are entrained in air flows which are jetted from a plurality of pulverizing nozzles, respectively, so as to collide against each other, and the powder material along with the air flow collide against each other, and then are pulverized. The air flow entirely fluidizes the powder material in the pulverization chamber, so as to accelerate pulverization caused by collision between the powder materials. Part of the powder material which has been pulverized and fluidized is guided to the area near a rotating rotor provided at the upper part of the pulverization chamber, and the particles of powder material each having a certain particle size or smaller are guided inside the rotor along with the fluid flow, and then powder as a final product (hereinafter referred to as product powder) is taken out from an outlet. The particles of the powder each having a certain particle size or larger are returned back to the outer periphery of the rotor by the centrifugal separation effect of the rotating roller, and are again returned back to the pulverization chamber, and then subjected to pulverization therein.\nFIG. 1 shows a cross-sectional view of a conventional fluidized bed pulverizing apparatus. With reference to FIG. 1, a structure of the conventional fluidized bed pulverizing apparatus and a method for producing powder will be described below. In FIG. 1, 1 denotes a powder material supply inlet, from which a powder material is supplied, 2 denotes an outlet which discharges pulverized powder as a product along with exhaust air, 3 denotes a centrifugal classification rotor which classifies the pulverized powder, 4 denotes a pulverization chamber in a fluidized bed container, 5 denotes fluid jetting nozzles whose jetting openings are arranged inside the pulverization chamber 4, and which face each other and jet fluid, 6 denotes a motor driving the centrifugal classification rotor 3. The external appearance of the main body of the entire fluidized bed pulverizing apparatus is a substantially cylindrical housing.\nThe operation of the fluidized bed pulverizing apparatus shown in FIG. 1 is as follows. At first, before operation of the apparatus, inside the pulverization chamber 4 a certain amount of the powder material is charged. Next, compressed air is jetted from each of the two fluid jetting nozzles 5 facing each other, the air jetted from each of the two fluid jetting nozzles 5 forms jetted air flow. The jetted air flow entrains the powder material which is present in the pulverization chamber 4, so as to transport the powder material. The two jetted air flows entraining the powder material collide against each other near the center of the pulverization chamber 4, so as to form air flow upward, downward, leftward and rightward directions inside the pulverization chamber 4. These air flows further entrain the powder material in the pulverization chamber 4, so as to form a fluidized bed of the powder material in the pulverization chamber 4. On the other hand, the powder material entrained in the jetted air flows collide against each other along with the collision of a plurality of jetted air flows, and are pulverized. Further, in the fluidized bed, the collision and pulverization of the powder material are repeated.\nAir in the pulverization chamber 4 passes from the outer periphery of the centrifugal classification rotor 3 located at the upper part of the pulverization chamber 4, through a gap between the rotors provided in the centrifugal classification rotor 3, and is guided to the outlet 2 connected to the centrifugal classification rotor 3, and then discharged from the outlet 2 to the outside. The powder material forming the fluidized bed is raised along with the exhaust air to the upper part inside the pulverization chamber 4, and enter the gap between the rotors from near the outer periphery of the centrifugal classification rotor 3. The centrifugal classification rotor 3 rotates at a certain rotating speed, among the powder material along with the air flow which reach the gap between the rotors, the powder material each having a certain particle size or larger is blown away to the outside of the centrifugal classification rotor 3 by centrifugal force. The particles of powder material each having a particle size smaller than a certain particle size along with the air flow are guided from the centrifugal classification rotor 3 to the outlet 2, and then discharged to the outside. The particles of powder material each having a certain particle size or larger are blown away to the outside of the centrifugal classification rotor 3, fall down in the pulverization chamber 4, and then are pulverized again in the fluidized bed.\nFrom the powder material supply inlet 1, the powder material in an amount corresponding to the amount of powder discharged from the outlet 2 are supplied to the pulverization chamber 4, and the amount of the powder material in the pulverization chamber 4 is kept constant. Thus, in the fluidized bed pulverizing apparatus, the particles of powder material each having a desired particle size can be continuously produced. Meanwhile, the particle sizes of the particles of the powder material discharged from the outlet 2 can be controlled by adjusting the rotating speed of the centrifugal classification rotor 3. The pulverizing speed of the powder material, namely production speed of the pulverized powder material can be controlled by adjustment of the speed and flow rate of the air flow jetted from the fluid jetting nozzle 5.\nIn the fluidized bed pulverizing apparatus, the powder material is repeatedly pulverized in the pulverization chamber, in order to obtain the particles of product powder each having a desired particle size. In this case, when the production speed of the product powder is intended to increase, it is necessary to increase air flow rate jetted from the fluid jetting nozzle 5 so as to increase the pulverization efficiency of the powder material. However, in the case where the air flow rate jetted from the fluid jetting nozzle 5 is increased, the amount of exhaust air is increased, decreasing the classification efficiency of the centrifugal classification rotor 3. As a result, the average particle size of the product powder may become large, or the particles of powder material each having a large particle size may be easily mixed in the product powder. The average particle size of the particles of product powder can be controlled by adjustment of the rotating speed of the centrifugal classification rotor 3 to some degree. However, it is not easy to prevent the large size particles of the powder material from being mixed in the product powder. Therefore, as a countermeasure for the problem of the mixture of the large size particles of the powder material in the product powder, there has been known a method of providing a baffle plate at the upper part of the pulverization chamber 4, by which the course particles are prevented from mixing in the product powder. However, this method may decrease pulverization efficiency, probably causing decrease in production speed.\nMoreover, there has been proposed a fluidized bed pulverizing apparatus (also referred to as “an air flow pulverizing apparatus”), for the purpose of improvement of pulverization efficiency of the fluidized bed pulverizing apparatus, adjustment of particle size of the product powder, and stabilization of product quality.\nFor example, Japanese Patent Application Publication (JP-B) No. 07-4557 discloses an air flow pulverization method, in which the pulverization efficiency of the powder material is improved by using a pulverization medium having a relatively large particle size.\nJapanese Patent Application Laid-Open (JP-A) No. 2002-126560 discloses an air flow pulverizer, in which the pulverization efficiency is improved by adjusting the internal pressure of a pulverization chamber to negative pressure, or rising temperature in the pulverization chamber.\nJapanese Patent (JP-B) No. 4025179 discloses an air flow pulverizer, in which a secondary collision unit for powder material which has collided by jetted air flow is provided so as to increase probability of collision between the powder materials, thereby increasing the pulverization efficiency.\nJP-B No. 4291685 discloses an air flow pulverizer, in which compressed air jetted from a jetting nozzle is heated so as to improve the pulverization efficiency of the powder material, and the particle size of the product powder is optimized.\nJP-A No. 2006-297305 discloses an air flow pulverizer, in which a space blocking member is provided in the inner wall of the pulverization chamber, particularly around a jetting nozzle, so as to decrease a dead space in the fluidized bed during the formation of the fluidized bed, thereby increasing the pulverization efficiency.\nJP-B No. 2503826 discloses an air flow pulverizing method, in which a bypass directly leading from the pulverization chamber to the channel for discharging the final powder is provided so as to control the particle size distribution of the product powder.\nJP-A No. 05-146704 discloses an air flow pulverizing method, in which a load current value of a motor for driving a classification rotor of a classifier is calculated as an integrated value of a predetermined time, and based on the value the supply amount of the powder material is adjusted so as to stabilize the particle size of the product powder.\nJP-A No. 3995335 discloses an air flow pulverizer which controls the quality of the product powder in such a manner that the density of fluidized powder material in the pulverization chamber of the air flow pulverizer and the amount of powder material deposited in the lower part of the pulverization chamber are measured, and according to the density and the amount, the taking out of the deposited powder material and the supply of the raw material of the powder are controlled.\nBy using the above-described fluidized bed pulverizing apparatuses (air flow pulverizing apparatuses) or the fluidized bed pulverization methods, a certain effect is obtained for the purpose of improvement of pulverization efficiency, adjustment of product quality, product quality stabilization. However, any of the fluidized bed pulverizing apparatuses (air flow pulverizing apparatuses) and the fluidized bed pulverization methods aims to improve pulverization efficiency, and adjust and stabilize of product quality only during steady operation. Therefore, they still have problems in terms of the adjustment of product quality and the stabilization of quality during the initial operation of the fluidized bed pulverizing apparatus.\nAt the beginning of the operation of the fluidized bed pulverizing apparatus, a powder material is entirely in non-pulverized state. When air is jetted from a jetting nozzle, the powder material present in the pulverization chamber is whirled up in the air jetted from the jetting nozzle, and the powder material is started to collide and form a fluidized bed. In the unsteady state during the initial formation of the fluidized bed, not only the abundance ratio of the pulverized particles of powder material each having a certain particle size or smaller is low, but also the ratio of the non-pulverized particles of the powder material each having a large particle size introduced into a centrifugal classification rotor provided at the upper part of the pulverization chamber is high. In such unsteady operation state, the particle size of the product powder discharged from the outlet along with the exhaust air from the centrifugal classification rotor tends to be large. Thus, during the initial operation of the apparatus, the quality of product powder is not stable. In the case where the quality of the product powder is emphasized, the discharged product powder is discarded or recycled as an off-specification product for a certain time during the initial operation of the fluidized bed pulverizing apparatus, until the quality of the product powder finally becomes stable. In the fluidized bed pulverizing apparatus which is expected to produce a large number of product lots, every time when operation restarts for changing a product lot, off-specification products are formed, causing significant decrease in the production efficiency."}
{"text": "The chemical, materials handling, and transportation fields preferably require containers having only a minimal number of ports and other openings thereon for safety, cost and other reasons. It is therefore advantangeous to provide containers such as cargo and storage tanks and bulk shipping containers having single openings which can be used for multiple purposes such as for filling and cleaning, for venting or relieving normally occurring pressure and vacuum conditions, and also for preventing dangerous pressure build-up in the container in the event of fire and the like. Government regulations can also require tanks, vessels and other containers for the transportation and storage of chemicals and other substances to have means providing prescribed flow rates for pressure and vacuum relief under predetermined pressure conditions, and also means providing temperature actuated pressure relief and depressurization at or above predetermined temperature conditions. Such relief means should be rugged and durable so as to be able to withstand abuse in industrial environments, should operate under severe adverse conditions, and should also be leak-resistant even in the event the container tips or overturns.\nThe device disclosed in U.S. Pat. No. 5,165,445, referenced above, mounts on or over a single opening or orifice on a container and can provide both pressure relief and vacuum relief for the container through the container opening. Such device, however, provides no means for temperature actuated pressure relief. Furthermore, such device does not include mounting means adapted for replacing container caps and covers and so does not readily facilitate using the same port or opening for alternative purposes such as for filling and cleaning the container.\nNumerous other devices attempting to satisfy one or more of the above discussed requirements are also known in the art. Reference is made to U.S. Pat. No. 5,031,790, which discloses a removable vented cap for fuel tanks including a self-actuating valve for pressure relief, an open vent for vacuum relief, and fusible plugs for temperature actuated pressure relief. Reference is also made to U.S. Pat. No. 3,385,468, which discloses a safety vent device for mounting over the vent hole of a gasoline tank, which device has passages for pressure and vacuum relief venting and a releasable member held in place with a thin band of fusible alloy to provide emergency pressure relief at a particular temperature. Such prior art devices are limited, however, as they include open passages and fusible means which can leak and allow contamination and such devices may also be damaged if the tank is overturned. Such devices can also expose the interior of a tank and the contents thereof to toxic fusible substances such as a solder substance containing lead alloy or like substances.\nContrasted to the relatively limited prior art devices discussed above, the subject invention relates to a durable and leak-resistant relief vent apparatus which mounts as a removable cover or cap over an opening on a container and can provide pressure relief and vacuum relief for the container at precise predetermined pressure and vacuum conditions and rapid fusible pressure relief at and above a precise predetermined temperature."}
{"text": "1. Field of the Invention\nAspects of the present invention relate to a fluorine-containing compound and an organic light-emitting device employing the same. More particularly, aspects of the present invention relate to a fluorine-containing compound that has electrical stability, good charge transport capability, and a high glass transition temperature and that can prevent crystallization. Aspects of the present invention further relate to an organic light-emitting device employing an organic layer including the fluorine-containing compound.\n2. Description of the Related Art\nElectroluminescent (EL) devices are self-emitting devices that have advantages such as a wide viewing angle, good contrast, and a rapid response time. EL devices are classified into inorganic EL devices, which include an emitting layer formed of an inorganic compound, and organic EL devices, which include an emitting layer formed of an organic compound. Organic EL devices show better brightness, driving voltage, and response speed characteristics compared to inorganic EL devices, and can create polychromatic light. Thus, extensive research into organic EL devices has been conducted.\nGenerally, organic light-emitting devices have a stacked structure including an anode, an organic light-emitting layer, and a cathode. A hole injection/transport layer or an electron injection layer may be further disposed between the anode and the organic light-emitting layer or between the organic light-emitting layer and the cathode to form an anode/hole transport layer/organic light-emitting layer/cathode structure, an anode/hole transport layer/organic light-emitting layer/electron transport layer/cathode structure, or the like.\nA polyphenyl hydrocarbon and an anthracene derivative have been described as material for forming a hole transport layer (U.S. Pat. Nos. 6,596,415 and 6,465,115).\nHowever, organic light-emitting devices including hole transport layers formed of currently available hole transport layer materials have disadvantages in terms of lifetime, efficiency, and power consumption characteristics, and thus, there is room for improvement in conventional organic EL devices."}
{"text": "Radiation is extensively used for the treatment of cancer and other diseased cells and tissues. Radiation therapy consists of exposing part or all of the body to a field of ionizing electromagnetic radiation. Often performed at 1 MeV or higher, the goal is to damage diseased cells. Although healthy cells frequently receive high radiation doses during such treatment, the healthy cells, ideally, are better able to repair the damage and remain viable while the diseased cells die.\nThe effectiveness of conventional radiation therapy is limited by insufficient radiation dosing due to the need to reduce radiation to normal cells and tissues. In many cases, the radiation dose to a tumor is the same as the dose to other tissues, especially surrounding tissues. This leads to significant toxicity in healthy cells. In order to increase the ratio of the dose to the intended target versus normal tissue (non target), radiation is often introduced to the tumor from different angles to reduce injury to skin and overlying tissue. However the x-rays also spread beyond the tumor and overshoot the target. The result is significant toxicity to an organism due to dosing of normal tissues.\nContrast agents are used to enhance the effect of x-rays for treatment of aberrant tissue. (U.S. Pat. Nos. 6,125,295 and 6,366,801). For example, a contrast agent is normally delivered to a tumor mass prior to delivering the radiation dose. These contrast agents have, as a component, an element with a high atomic number (Z), such as iodine or gold. The interaction between the ionizing radiation and the greater cross-section of the high Z material creates additional ionizations that result in greater cell toxicity at the site of the tumor.\nContrast agents also improve the accuracy of assessing a disease state. To be useful, the contrast material must be delivered to the area where a suspected abnormality may be present for radiation exposure to result in high enough contrast for a successful diagnosis\nHowever, conventional contrast agents have the disadvantage that they lack affinity for the cells and tissues to be treated so that the residence time of the agent in the targeted tissue is short. The poor uptake of the conventional contrast agent by a tumor means that the agent needs to be applied directly into the tumor. Furthermore, the contrast agents migrate out of the tumor quickly and delivery of the radiation is required very soon after administering the contrast agent, often within one hour. If delivered by intravenous administration, common contrast agents often require relatively large volumes of contrast agent solution to be administered within a short period of time e.g., 100 ml within one minute. This creates a risk of rapid allergic reaction and can cause discomfort to the patient.\nTo achieve more specific cellular and tissue targeting with these agents, they are typically modified using a biological carrier such as a protein, or a monoclonal antibody, or fragments thereof. Thus a monoclonal antibody combined with a payload of iodine or other heavy element can be used to more selectively deliver high Z atoms to a tumor. These agents have been shown to be useful for the treatment (when the elements themselves are radioactive isotopes) and diagnosis of cancers. However, the biodistribution of these systems is unfavorable to enhance radiation therapy when the elements are not radioactive. In addition, the retention time of the dose in blood when radioactive elements are used is long, and usually only a small portion of the dose is observed at the site of the tumor. Unfortunately, the density of antigen sites that the tumor can present is low and so the potential amount of high Z material that can be delivered is relatively low. Better agents are needed to enhance the effect of radiation at tumor sites."}
{"text": "This invention relates to a picture frame and a picture frame assembly, in particular, a one-piece picture frame made from a suitably-foldable material for framing material to be framed wherein the material to be framed comprises at least one flat, planar object.\nIn the past, there have been many foldable frames made from foldable material. However, many of those frames had pictures pre-printed on the front surface of the back piece of the frame, or had the object to be framed adhered to the front surface of the back piece of the picture frame.\nMoreover, the prior art picture frames were either too simple in that they did not provide a suitable frame into which the material to be framed could be inserted. Rather, the material to be framed was merely placed on the frame and attached to the frame. Such frames were not suitable for framing combined material comprising the object to be framed, such as a picture, photograph or postcard, plus a suitable viewing sheet over the object to be framed through which the object to be framed could be viewed, a suitable backing material and perhaps even matting.\nOn the other hand, other of the prior art picture frames were too complicated to be easily and cheaply manufactured and to be readily used. This is because many of these prior art picture frames utilized a plurality of intricate tabs and tab receiving areas. As a consequence, the picture frame required too many cuts, folds and insertions.\nAlso, none of the prior art provided a suitable frame incorporating as part of the frame a supporting means to support the frame on, for example, a desk top."}
{"text": "This section is intended to provide a background or context to the invention that is, inter alia, recited in the claims. The description herein may include concepts that could be pursued, but are not necessarily ones that have been previously conceived or pursued. Therefore, unless otherwise indicated herein, what is described in this section is not prior art to the description and claims in this application and is not admitted to be prior art by inclusion in this section. Thin film materials of metal-metal oxides of nanocomposites can have many applications, including resistive layers for electronic applications, such as, for example, electron multipliers like microchannel plates, resistive memories, electro-chromic devices, biomedical devices and charge dissipating coatings on micro-electromechanical systems. A related application, U.S. Ser. No. 13/011,645, which is incorporated by reference herein, describes microchannel plate fabrication by atomic layer deposition (“ALD” hereinafter), which provide an example of how one can benefit from the tunable resistance coatings and methods of preparation described herein."}
{"text": "The shipping industry is comprised of various types of shippers, including small package carriers, less-than-a-load (“LTL”) carriers, and truck load carriers. Small package carriers usually transport packages or boxes from multiple consignors, while truck load carriers typically transport entire trailer loads from a single consignor. LTL carriers, on the other hand, generally transport freight that falls in between small packages and trailer loads. For example, LTL carriers may move freight from multiple consignees in a single trailer load, such as crates, scrap metal banded together, vehicle parts, pallets of boxes, drums, and the like. This freight is usually consolidated into a single trailer and transported through a carrier's transportation network. To track the freight and provide carrier personnel with routing and handling instructions, LTL carriers currently use paper bills of lading. The use of paper bills of lading decreases the efficiency and throughput of carrier transportation networks—relying almost solely on the efficiency of carrier employees—and does not provide for real-time visibility of freight progressing through a transportation network.\nIn addition to using paper bills of lading to track freight, LTL carriers often use the weight and shipping classifications on the bills of lading (provided by consignors) to appropriately charge consignors and/or consignees for transporting the freight. And although many carriers have internal audit mechanisms to verify the weight and shipping classifications provided on the bills of lading, the audit procedures are generally paper-driven and manual in nature. The paper-driven and manual nature of the current audit procedures limit the carriers' ability to efficiently and cost-effectively audit much of the freight they transport. For instance, carriers often have personnel who are specifically employed to audit freight shipments. Typically, the personnel manually verify the shipping classifications and weight provided on the bills of lading by personally inspecting the freight. This methodology usually enables carriers to only audit select freight shipments, leaving the majority of freight shipments unaudited and possibly incorrectly classified."}
{"text": "This invention relates to a food cutting apparatus and particularly to apparatus for dividing a block of cheese into a plurality of relatively small chunks suitable for grating.\nThe mass marketing of various food products has resulted in development of automatic food forming apparatus. For example, pizzas which are mass marketed may advantageously be constructed with various automatic machinery. A highly satisfactory automated system is disclosed in Applicant's recently issued U.S. Pat. No. 3,779,205, which generally discloses an in-line processing line including means for delivering blocks of cheese to an automatic slicer which divides the block into a plurality of smaller pieces which are fed to a power-driven shredder. The shredded cheese is fed to an automatic depositing apparatus for depositing of the grated cheese upon the crust or dough bases in combination with automatic apparatus for depositing of meat, sauce and the like. Such a structure may provide a high speed, mass production of pizzas which are frozen and then shipped to various marketing outlets. The automated forming not only improves the quality of the product but significantly reduces the per unit cost.\nIn all of such automatic apparatus, very stringent requirements are established with respect to the construction of the apparatus to ensure high degrees of sanitary conditions. Generally, the apparatus will be made of stainless steel and must be made to permit convenient, complete cleaning of the apparatus.\nThe cheese cutting and grating or shredding means must correspondingly be constructed to permit complete and reliable cleaning. A practical, automated apparatus has not been commercially available which has a large capacity and a convenient construction suitable for high level sanitary application."}
{"text": "The present invention relates generally to lockers and more particularly to locker door retrofit assemblies. The present invention is designed to replace existing metallic door or other assemblies with a locker door assembly comprised of a more suitable material.\nLockers for storing clothing, articles of merchandise, etc., are commonly constructed. Storage lockers are found in many different settings. For example, athletic facilities have lockers to allow athletes to store their possessions while participating in athletic events. Community swimming pools typically have lockers for storing street clothes while a person swims. Lockers are also found in industry where they are used for several purposes, such as the storage of equipment, work clothes, chemicals, and other items which are preferably kept in such a concealed environment when not in use. Lockers are also commonly found in airports, hospitals, school hallways, prisons, and many other sites too numerous to mention.\nMost commonly steel sheet metal is used as a primary construction material with metal fasteners used to assemble the finished locker. Metallic lockers suffer from several disadvantages. They are easily damaged or marred in some way such as by dents, scratches and graffiti. Moreover, the metal is subject to damage from rust, odors, delamination and fading.\nAttempts have been made to solve the above-mentioned problems by wholly replacing these metallic lockers with plastic or wood lockers. Some storage lockers were built into either an existing wall of a building or into a building wall while under construction. These in-wall lockers may be expensive to replace with plastic lockers. These problems prevent many from taking advantage of the properties that plastic offers over steel sheet metal.\nA need exists for a locker that is designed to overcome the aforementioned disadvantages. The present invention is a locker door retrofit assembly comprised of a material designed to overcome the above disadvantages, such as a plastic or a composite material. The following disclosure describes a plastic locker door retrofit assembly. However, it must be understood that any non-metal material that exhibits the desired characteristics may be utilized for the present invention.\nIn the present invention existing metallic doors and metallic jams or doors of other materials, such as wood, may be replaced with the plastic locker door assembly. Due to the plastic construction of the door assembly, the locker face will be resistant to many forms of abuse that lockers commonly receive. The locker doors will not dent as will metal lockers. The locker door of the present invention will maintain its color throughout its entire cross-section. Due to the preferred homogeneous nature of the plastic door assembly, the lockers of the-present invention will not delaminate. Furthermore, most materials used in the application of graffiti are readily removed from plastic panels to return the locker doors to the original surface appearance. The remaining metallic body of the locker system is hidden behind the plastic face and is thus protected. Moreover, the substitution with a plastic door assembly provides many cosmetic and aesthetically pleasing attributes to the locker system. These plastic lockers may carry almost any color scheme desirable. Colors may be chosen to match the surrounding decor, to provide a color coding scheme and/or to provide a medium for an organizational theme.\nOther features and advantages of the present invention will be apparent from the following description and claims and are illustrated in the accompanying drawings which show preferred features of the present invention and the principles thereof."}
{"text": "This invention generally relates to information retrieval, and more specifically, to assembling answers from multiple documents. Even more specifically, embodiments of the invention relate to Question Answering systems and methods implementing parallel analysis for providing answers to questions and in which candidate answers may be assembled from multiple documents.\nGenerally, question answering (QA) is a type of information retrieval. Given a collection of documents (such as the World Wide Web or a local collection), a QA system should be able to retrieve answers to questions posed in natural language. QA is regarded as requiring more complex natural language processing (NLP) techniques than other types of information retrieval such as document retrieval, and QA is sometimes regarded as the next step beyond search engines.\nQA research attempts to deal with a wide range of question types including: fact, list, definition, how, why, hypothetical, semantically-constrained, and cross-lingual questions. Search collections vary from small local document collections, to internal organization documents, to compiled newswire reports, to the world wide web.\nClosed-domain question answering deals with questions under a specific domain (for example, medicine or automotive maintenance), and can be seen as an easier task because NLP systems can exploit domain-specific knowledge frequently formalized in ontologies. Alternatively, closed-domain might refer to a situation where only a limited type of questions are accepted, such as questions asking for descriptive rather than procedural information. Open-domain question answering deals with questions about nearly everything, and can only rely on general ontologies and world knowledge. Open-domain Q/A systems, though, usually have much more data available from which to extract the answer.\nAccess to information is currently dominated by two paradigms: a database query that answers questions about what is in a collection of structured records; and a search that delivers a collection of document links in response to a query against a collection of unstructured data (text, html etc.).\nOne major challenge in such information query paradigms is to provide a computer program capable of answering factual questions based on information included in a large collection of documents (of all kinds, structured and unstructured). Such questions can range from broad such as “what are the risk of vitamin K deficiency” to narrow such as “when and where was Hillary Clinton's father born”.\nUser interaction with such a computer program could be either a single user-computer exchange or a multiple turn dialog between the user and the computer system. Such dialog can involve one or multiple modalities (text, voice, tactile, gesture etc.). Examples of such interaction include a situation where a cell phone user is asking a question using voice and is receiving an answer in a combination of voice, text and image (e.g. a map with a textual overlay and spoken (computer generated) explanation. Another example would be a user interacting with a video game and dismissing or accepting an answer using machine recognizable gestures or the computer generating tactile output to direct the user.\nThe challenge in building such a computer system is to understand the query, to find appropriate documents that might contain the answer, and to extract the correct answer to be delivered to the user. Currently, understanding the query is an open problem because computers do not have human ability to understand natural language nor do they have common sense to choose from many possible interpretations that current (very elementary) natural language understanding systems can produce.\nBeing able to answer a factual query in one or multiple dialog turns is of great potential value as it enables real time access to accurate information. For instance, advancing the state of the art in question answering has substantial business value, since it provides a real time view of the business, its competitors, economic conditions, etc. Even if QA is in a most elementary form, it can improve productivity of information workers by orders of magnitude.\nU.S. patent application Ser. No. 12/152,441, the disclosure of which is hereby incorporated herein by reference in its entirety, describes a QA system involving the generation of candidate answers and selecting a final answer (or ranking a list of final answers) from among the set of candidate answers.\nCurrent information retrieval and question answering systems attempt to satisfy a user's information need by identifying the single document segment (e.g., entire document, contiguous sequence of one or more sentences, or a single phrase) that is most likely to contain relevant information. There are many information needs, however, that cannot be satisfied by a single document segment."}
{"text": "Breastmilk pumps are well known and are generally comprised of a hood that fits over the breast, a vacuum pump connected to the hood for generating an intermittent vacuum within the hood, and a receptacle for the expressed milk. Manually driven vacuum pumps as well as those that are driven by a motor are ordinarily used. The vacuum pumps of these devices, as a rule, intermittently generate a vacuum or negative pressure within the hood, with the hood encompassing the nipple and a substantial amount of the breast. The intermittent suction action of the pump serves to pull on the breast and thereby extract milk in an action reminiscent of suckling. The milk so extracted typically flows from the hood into a collection container, e.g., a bottle, for storage and later use. A breastpump of the foregoing type is shown in U.S. Ser. No. 07/053,055, filed May 22, 1987, now U.S. Pat. No. 4,857,051.\nApart from the purely hygienic requirements for such equipment, there are also certain technical problems to consider. One such significant problem is that varying degrees of vacuum can be generated as the milk receptacle fills, which must then be compensated for. A solution to this problem is to provide a valving mechanism which serves to regulate the negative pressure applied."}
{"text": "1. Field of the Invention\nThis invention relates to EEPROM devices and more particularly to bit erasing as well as block erasing therein.\n2. Description of Related Art\nIn the past in EEPROM devices, erasure has been performed groups of eight bits known as a byte, byte by byte, but one cell needs 2 1/8 transistors and tunneling window, so the cell size is large. See FIG. 1B where a cell includes a cell transistor Tc, a bit select transistor T.sub.bit, and share the byte select transistor T.sub.byte with seven other cells.\nAn alternative to the EEPROM which solves part of the problem of large cell size is the flash EPROM. While the flash EPROM has the advantage that the cell size is small, the problem is that the flash EPROM has the disadvantage that it erases block by block, and one can not use it to erase bit by bit or byte by byte.\nFIG. 1A shows a cross section of a prior art device known as a FLOTOX (Floating-Gate Tunneling Oxide) EEPROM cell known as an E.sup.2 PROM cell which requires a tunnel oxide window and a select transistor, not shown in FIG. 1A, so the cell size is large. The FLOTOX cell is described in Samachisa et al of SEEQ Technology Inc. and U of C, Berkeley for \"A 128k Flash EEPROM Using Double-Polysilicon Technology\", IEEE J. Solid-State Circuits Vol. 8C-22, No. 5, pp 676-683 (Oct. 1987.) In FIG. 1A the device includes a P-substrate 10 containing a drain region 11 connected to V.sub.d voltage source, a source region 12 connected to V.sub.s voltage source, a tunnel oxide 14, gate oxide layer 15, a field oxide (FOX) 16, a floating gate 17 composed of polysilicon 1 first dielectric layer 18, a control gate 19 composed of polysilicon 2 connected to V.sub.g voltage source, and a second dielectric layer 20. The tunnel oxide 14 is located between the floating gate 17 and the N+ drain region 12.\nFIG. 1B shows a prior art E.sup.2 PROM cell array (in accordance with FIG. 1A) and shows the connection lines required for its operation. Referring to Table 1 an operation table shows how an E.sup.2 PROM can program and erase in groups of cells, by the byte (8 bits.)\nTABLE 1 ______________________________________ PROGRAM ERASE READ ______________________________________ SELECTED WORD LINE 20 V 20 V 5 V UNSELECTED WORD 0 V 0 V 0 V LINE PROGRAM LINE 17 V 0 V 0 V BIT LINE 0 (ERASED) 0 V 17 V 1.6 V BIT LINE 7 0 V 0 V 2.0 V (PROGRAMMED) ______________________________________\nFIG. 2A shows a prior art flash memory cell structure know as an ETOX.RTM. flash memory comprising a flash memory cell with tunnel oxide TO below the floating gate FG. The device includes source S and drain D in the substrate and a control gate CG above the floating gate FG separated therefrom by a dielectric layer, with the erase E function from floating gate FG to source indicted and the programing P function from drain to floating gate FG indicated.\nFIG. 2B shows a cell array of the prior art device of FIG. 2A. The array is controlled by bit lines BL1, BL2, BL3 and BL4 which extend vertically to the cells. The bit lines are connected to the S/D circuits of the cells, which are connected at the opposite ends to one of a plurality of voltage supply sources represented by V.sub.ss1, V.sub.ss2, and V.sub.ss3. Word lines WL1, WL2, and WL3 are the horizontally directed lines connected to the cells. In particular, the word lines WL1, WL2, and WL3 are connected to the control gates of the cells. FIG. 2B along with TABLE 2 which is a flash memory operation table shows how such an ETOX flash memory device operates. In particular, Table 2 below shows the operation table for the ETOX flash memory of FIGS. 2A and 2B is unable to program and erase by byte (8 bits.) Referring to Table 2, one can program or read a single bit (cell,) but one can not erase a single bit, because the word lines WL1, WL2, . . . WLn and V.sub.ss1, V.sub.ss2. . . V.sub.ssn lines are parallel with each other (i.e. are oriented in the same direction. )\nTABLE 2 ______________________________________ PROGRAM ERASE READ ______________________________________ SELECTED WORD LINE 12 V 0 V 5 V (WL2) UNSELECTED WORD 0 V 0 V 0 V LINE (WL1, WL3) SELECTED BIT LINE (2) 7 V 0 V 1.6 V UNSELECTED BIT 0 V 0 V 0 V LINE 1, 3, 4 SELECTED V.sub.ss (2) 0 V 12 V 0 V UNSELECTED V.sub.ss (1, 3) 0 V 0 V 0 V ______________________________________"}
{"text": "Coordinated multi-point (CoMP) is an interference avoidance concept that can be used to improve system spectral efficiency and cell edge user throughput performance. CoMP may be used to avoid interference to other cells by coordination of the transmissions across multiple eNBs."}
{"text": "For the production of perchlorate, various anode materials have been used commercially, including smooth massive platinum, platinized titanium or tantalum (despite a tendency to produce excess oxygen) and lead dioxide coated on titanium or graphite, although these lead dioxide anodes have a high overvoltage and wear rapidly.\nSome proposals have already been made to combine platinum group metals and tin dioxide in electrode coating materials. For example, U.S. Pat. No. 3,701,724 mentioned an anode for chlorine production having a coating consisting essentially of a minor amount of a platinum group metal and/or platinum group metal oxides with a major amount of SnO.sub.2, Sb.sub.2 O.sub.5, Sb.sub.2 O.sub.3 or GeO.sub.2 and mixtures thereof. However, the claims and examples of this patent are directed solely to such coatings containing platinum group metal oxides and there is no enabling disclosure of a coating containing a platinum group metal. Also, U.S. Pat. No. 3,882,002 proposed an anode for chlorine production having a valve metal substrate coated with an intermediate layer of tin dioxide which was covered with an outer layer of a platinum group metal or oxide thereof. Neither of these proposals was directed to improving electrolytic performance in the production of percompounds."}
{"text": "This invention relates to a surgical implant adapted for secure attachment to living bone. More particularly this invention relates to an improved orthopedic or dental implant formed to have bone contacting surfaces comprised of biocompatible organic polymers substituted with oxyacid groups or salts thereof. Applicant has found that the presence of such oxyacid substituted polymers, at least at the bone contacting surface, promotes interfacial osteogenesis and provides sites to which bone can chemically bond, thus fostering direct chemical bonding of the implant surface with the biological polymers present in developing bone tissue.\nIn severe cases of arthritis or other bone and joint degenerative diseases, surgical replacement of the affected joint and bone tissue is a commonly used procedure. In cases of dentition lost to disease or trauma, endosseous dental implants have been used to restore function. Research and development efforts have been highly successful in identifying materials for use alone or as composite structures for prosthetic implants --materials which meet both the basic physical (biomechanical) demands and the chemical biocompatibility requirements dictated by their use in implanted devices. Due to the stringent mechanical demands placed on load bearing bone prostheses, metals have been the material of choice for the most severely loaded parts of such implants. The metals generally used in load bearing components of orthopedic devices are limited to the cobalt, chromium, molybdenum alloys, titanium and surgical stainless steel. More recently, high strength ceramics and reinforced polymers have been introduced for such applications.\nNotwithstanding the major advances which have been made in implant materials development, patients still face the trauma and expense of implant failure. The most common point of bone implant failure is not breakage or failure of the prosthetic implant itself; the more common failure of implants is at the living bone-implant interface. In other words, the implant simply works loose from its implanted position. Early total joint replacements were fixed to the bone of a recipient through a press fit of the prosthesis into a carefully prepared surgical bed. This method often resulted in loosening of the implant in the long term.\nA review of the recent literature reveals significant research and development efforts directed to improving implant-bone fixation. A major advance in joint replacement surgery was the introduction of the use of poly(methyl methacrylate) [PMMA] bone cement for fixing the components of a joint prosthesis to bone. PMMA is not a glue or adhesive but a true cement which works mechanically. It is applied in a dough-like state as a grouting agent between the bone and the implant so that it can flow around the contours of the bone and the implant and into the interstices of cancellous bone. Upon hardening it forms a mechanically interlocked attachment between the bone and the implant. While PMMA bone cement provides a secure fixation of the prosthesis with living bone in the short term, the long term loss of implant fixation has proven to be a significant problem. The degeneration in implant fixation begins with a resorption of the bony tissue immediately adjacent to the bone cement and the replacement of that tissue with a soft fibrous tissue capsule. Since the fibrous tissue is far more compliant than bone, the thicker the capsule the looser the implant becomes. Since the thickness of the capsule tends to increase with motion the loosening process is self-reinforcing.\nIn an alternative method of implant fixation enjoying widespread use, the implant is provided with a highly porous surface coating that provides interstices into which bone can grow. Materials which provide pore size distributions of 50 to 500 microns have shown considerable promise and have become more widely used especially in young, active patients. For bone to interlock with the pore structure of the implant the implant must be firmly fixed at the time of surgery and load application must be minimized during the in growth period. Immediate surgical fixation is usually accomplished by the mechanical impaction of the implant into a slightly sub-sized surgical bed. The problems related to both of these fixation methods arise from the lack of affinity demonstrated by healing bone for the heretofore known metallic alloys and polymeric materials used in reconstructive orthopedic and oral surgery. When these materials are placed into bony defects, bone does not deposit directly onto the implant surfaces but, significantly, remains separated from them by at least a thin layer of soft tissue. This precludes the possibility of any chemical bond between the bone and the implant and limits the fixation modes to those based on mechanical interlock. Materials of this type are defined as \"osteophobic\" for the purposes of this disclosure.\nDue to the problem of obtaining fixation to living bone by purely mechanical means, more recent research efforts have been directed at finding more bone-tissue-compatible implant materials, particularly materials having some demonstrable chemical affinity for regenerating bone tissue. It has been found that when certain minerals are implarted into osseous defects, newly forming bone will deposit directly onto the mineral surface without any intervening layer of soft tissue. Furthermore not only does bone deposit directly on the surface of these minerals, but it has been observed that bone will adhere even to a smooth surface of said minerals through the formation of chemical bonds across the bone-mineral implant surface.\nFor the purposes of this disclosure materials onto which bone deposits and chemically bonds are defined as \"osteophilic\", and the process whereby bone deposits directly onto a free surface is defined as \"interfacial osteogenesis.\" The first materials recognized to have osteophilic properties were hydroxyapatite [Ca.sub.10 (PO.sub.4).sub.6 (OH).sub.2 ] and tricalcium phospate [Ca.sub.3 (PO.sub.4).sub.2 ]. Their osteophilicity has been rationalized by their close similarity in chemical structure with the mineral phase of mammalian bone. More recently, this inventor has made the surprising discovery that chemical similarity with bone structure is not necessary for a mineral substance to exhibit osteophilic properties. Indeed, it has been found that calcite [CaCO.sub.3 ], dolomite [CaMg(CO.sub.3).sub.2 ] (See my copending U.S. patent application Ser. No. 251,225), and more recently the minerals fluorapatite, aragonite, magnesite, witherite and barite have been found to exhibit osteophilic properties. A number of researchers have sought to take advantage of in vivo chemical bonding of bone to bioactive mineral surfaces. See, for example, U.S. Pat. Nos. 3,787,900; 3,919,723; 4,168,326; 4,320,514; 4,366,188; 4,373,217; and 4,437,192.\nThe present invention is based on the discovery that certain organic polymers, that is, organic polymers bearing salt-forming oxyacid functional groups, exhibit an affinity toward developing bone tissue similar to that which has been observed for the above-mentioned inorganic mineral materials. The brittle nature of osteophilic minerals limits their usefulness to low stress applications. One such application is as a coating for porous-surfaced dental and orthopedic implants. It is known that pore sizes on the order of about 100-200 microns diameter are required for bony in growth into osteophobic materials. However, bone will grow into cracks and pores as small as 2 microns in diameter in osteophilic minerals. The pores of osteophobic metals, for example, have been coated with a thin layer of an osteophilic mineral to promote the rate of bony in growth. One serious problem related to the application of thin mineral coatings to porous metal surfaces is the lack of adhesion of such coatings to the metal substrate and the propensity for the minerals to dissolve in vivo when applied as thin films.\nThe present invention is based on the discovery that bone will deposit onto and chemically bond with organic polymeric materials having, at least on their bone contacting surfaces, covalently attached salt forming oxyacid functional groups. That discovery coupled with the inherent versatility of organic polymers for implant fabrication represents a significant advance in the art. Polymeric materials exhibit superior processability and mechanical properties. They can be extruded, molded, shaped and machined, or they can be applied, as a melt or dissolved in solution, to coat the surfaces of prosthetic devices constructed of other materials or material composites. Many biocompatible polymers in bulk form exhibit weight/strength properties (especially when reinforced with high tensile strength fibers) unmatched by the metal or metal/ceramic materials which have been used for construction of prosthetic devices. Still a further advantage of the use of biocompatible organic polymers for construction of orthopedic protheses is their chemical versatility. That is, implant devices constructed of or coated with biocompatible organic polymers can be surface modified by chemical treatment using a wide variety of reaction conditions to bear the desired surface-active oxyacid groups in accordance with this invention.\nIt is therefore an object of this invention to provide a method for promoting interfacial osteogenesis on an implant surface.\nIt is another object of this invention to provide a method for promoting bone in growth and bone adhesion to bone contacting surfaces of implanted prostheses.\nAnother object of this invention is to provide orthopedic prostheses having improved bone contacting surfaces comprised of a biocompatible organic polymer substituted with salt-forming oxyacid groups selected from the group consisting of the oxyacids of carbon, sulfur and phosphorus."}
{"text": "I. Field of the Invention\nThe present invention relates generally to utility bags used to contain a variety of articles, and more particularly to such bags which includes means for holding the bag in an open configuration through the use of stiffening members.\nII. Prior Art\nOver the years, a wide variety of utility bags have been devised to carry virtually any type of article. Many of those bags are constructed from a flexible material, such as a fabric, cloth, leather or synthetic sheet, and some include means for keeping the bag in an open configuration.\nOf those bags which employ means for keeping the bag in an open configuration, this objective is often accomplished by the inclusion of curved rods sewn into the upper rim of the opening. Other such bags may include strips that resiliently bias the sides away from one another, such as the bag disclosed in U.S. Pat. No. 4,561,525. In that device, the sides are allowed to fold onto themselves when the bag is opened, while a pair of strips of heavier gauge flexile material attached to the inside of the bag act to keep the bag open.\nWhile the aforementioned bag may achieve its stated purpose, there remains a need for a utility bag which can easily and reliably be opened and held open by virtue of its own structure, and through the folding of its sides. Ideally, such a bag would allow folding away of the sides in such a manner as to expose the entire internal surface area of the base of the bag. Furthermore, the bag should allow an upper portion of the sides to conveniently fold parallel to the lower portion of the sides, such that the open configuration of the bag is maintained by the filling of the bag with the desired articles. In this manner, the user has easy access to the entire contents of the bag without having to push away the sides. As will be seen in the descriptions provided herein, such a bag is particularly useful for holding vertically oriented, relatively flat plastic bags, such as those designed to hold plastic worms used by fisherman. However, it will become apparent that the construction of the present invention and its corresponding unique features may also be used in a wide range of situations, such as a cosmetic bag, shaving kit, and other diverse applications."}
{"text": "Field effect transistors (FETs) are widely used in the electronics industry for switching, amplification, filtering, and other tasks related to both analog and digital electrical signals. Most common among these are metal oxide semiconductor field effect transistors (MOSFETs), wherein a gate electrode is energized to create an electric field in a channel region of a semiconductor body, by which electrons are allowed to travel through the channel between a source region and a drain region of the semiconductor body. The source and drain regions are typically formed by adding dopants to targeted regions on either side of the channel. A gate dielectric or gate oxide is formed over the channel, and a gate electrode or gate contact is formed over the gate dielectric. The gate dielectric and gate electrode layers are then patterned to form a gate structure overlying the channel region of the substrate.\nIn operation of the resulting MOS transistor, the threshold voltage (Vt) is the gate voltage value required to render the channel conductive by formation of an inversion layer at the surface of the semiconductor channel. Complimentary MOS (CMOS) devices have become widely used in the semiconductor industry, wherein both n-channel and p-channel (NMOS and PMOS) transistors are used to fabricate logic and other circuitry. For enhancement-mode (e.g., normally off) devices, the threshold voltage Vt is positive for NMOS and negative for PMOS transistors. The threshold voltage is dependent upon the flat-band voltage, where the flat-band voltage depends on the work function difference between the gate and the substrate materials, as well as on surface charge.\nThe work function of a material is a measure of the energy required to move an electron in the material outside of a material atom from the Fermi level, and is usually expressed in electron volts (eV). For CMOS products, it is desirable to provide predictable, repeatable, and stable threshold voltages (Vt) for the NMOS and PMOS transistors. To establish Vt values, the work functions of the PMOS and NMOS gate contact and the corresponding channel materials are independently tuned or adjusted through gate and channel engineering, respectively.\nGate stack engineering is employed to adjust the work function of the gate contact materials, where different gate work function values are set for PMOS and NMOS gates. The need to independently adjust PMOS and NMOS gate work functions has made polysilicon attractive for use as a gate contact material in CMOS processes, since the work function of polysilicon can be easily raised or lowered by doping the polysilicon with p-type or n-type impurities, respectively. The PMOS polysilicon gates are typically doped with p-type impurities and NMOS gate polysilicon is doped with n-type dopants, typically during implantation of the respective source/drain regions following gate patterning. In this way, the final gate work functions are typically near the Si conduction band edge for NMOS and near the valence band edge for PMOS. The provision of dopants into the polysilicon also has the benefit of increasing the conductivity of the gate electrode. Polysilicon has thus far been widely used in the fabrication of CMOS devices, wherein the gate engineering provides a desired gate electrode conductivity (e.g., sheet resistance value) by conventional tuning (e.g., implants), and the threshold voltage fine tuning is achieved by tailoring the channel doping level through the Vt adjust implants.\nFIG. 1 illustrates a conventional CMOS fabrication process 10 beginning at 12, in which front end processing is performed at 14, including well formation and isolation processing. At 16 and 18, channel engineering is performed (e.g., Vt adjust, punch-thru, and channel stop implants) for PMOS and NMOS regions, respectively. A thin gate dielectric and an overlying polysilicon layer are formed at 20 and 22, respectively, and the polysilicon is patterned at 24 to form gate structures for the prospective NMOS and PMOS transistors. The gate structures are then encapsulated at 26, typically through oxidation, and highly-doped drain (HDD) implants are performed at 28 to provide p-type dopants to prospective source/drains of the PMOS regions and n-type dopants to source/drains of the NMOS regions, using the patterned gate structures and isolation structures as an implantation mask. Sidewall spacers are then formed at 30 along the lateral sidewalls of the gate structures.\nAt 32, the PMOS source/drain regions and the PMOS polysilicon gate structures are implanted with p-type dopants to further define the PMOS source/drains, and to render the PMOS gates conductive. Similarly, the NMOS source/drain regions and the NMOS polysilicon gate structures are implanted at 34 with n-type dopants, further defining the NMOS source/drains and rendering the NMOS gates conductive. Thereafter, the source/drains and gates are silicided at 36 and back end processing (e.g., interconnect metalization, etc.) is performed at 38, before the process 10 ends at 40. In the conventional process 10, the channel engineering implants at 16 and 18 shift the Vt of the PMOS and NMOS channel regions, respectively, to compensate for the changes in the PMOS and NMOS polysilicon gate work functions resulting from the source/drain implants at 32 and 34, respectively. In this manner, the desired work function difference between the gates and channels may be achieved for the resulting PMOS and NMOS transistors, and hence the desired threshold voltages.\nThe gate dielectric or gate oxide between the channel and the gate electrode is an insulator material, typically SiO2, nitrided SiO2, or other dielectric, that operates to prevent current from flowing from the gate electrode into the channel when a voltage is applied to the gate electrode. The gate dielectric also allows an applied gate voltage to establish an electric field in the channel region in a controllable manner. Continuing trends in semiconductor product manufacturing include reduction in electrical device feature sizes (scaling), as well as improvements in device performance in terms of device switching speed and power consumption. MOS transistor performance may be improved by reducing the distance between the source and the drain regions under the gate electrode of the device, known as the gate or channel length, and by reducing the thickness of the layer of gate dielectric that is formed over the semiconductor surface.\nHowever, there are electrical and physical limitations on the extent to which SiO2 gate dielectrics can be made more thin. These include gate leakage currents tunneling through the thin gate oxide, limitations on the ability to form very thin oxide films with uniform thickness, and the inability of very thin SiO2 gate dielectric layers to prevent dopant diffusion from the gate polysilicon into the underlying channel. Accordingly, recent scaling efforts have focused on high-k dielectric materials having dielectric constants greater than that of SiO2, which can be formed in a thicker layer than scaled SiO2, and yet which produce equivalent field effect performance. A thicker high-k dielectric layer can thus be formed to avoid or mitigate tunneling leakage currents, while still achieving the required electrical performance equivalent (e.g., capacitance value) to a thinner SiO2.\nIt has also been proposed to utilize hafnium-based high-k dielectric materials in combination with a lanthanide series metal to lower the work function of metal gates. The lanthanide series metal is provided as a distinct surface layer over the high-k dielectric material. This proposal, however, may decrease the overall equivalent oxide thickness (EOT) of the layer of gate oxide."}
{"text": "The use of rechargeable nickel-cadmium (nicad) batteries for consumer products is well established. Such rechargeable batteries are frequently used in portable power tools, such as cordless power drills, saws and the like. Additionally, rechargeable batteries also find application in shavers, photographic equipment and other products.\nUnlike disposable batteries, however, the nicad batteries require recharging upon dissipation of the electrical energy stored therein. The recharging period of the nicad batteries, if too long, may thus diminish the effectiveness of the power tools which incorporate the batteries. There have thus been prior art attempts to speed up the charge rate in order more quickly to restore the batteries to full capacity.\nThe normal recommended continuous charge rate for nicad batteries is C/10 where C is the battery capacity in ampere-hours. The normal charge rate thus results in a time of 12 hours or more to recharge a battery pack. Such a time requirement is excessive, however. If the batteries powering a product are discharged prior to completion of the desired task, it is necessary for the operator either to wait for a recharge or to replace the battery pack with a fully charged replacement pack. The first approach, as above noted, is typically highly time consumptive while the second is expensive.\nAccordingly, the prior art has developed several approaches to reducing the recharge time for rechargeable batteries, including various techniques to avoid overcharging the units.\nIn one approach to the problem, battery manufacturers have conducted research into battery characteristics under charge and have developed special cells. Thus, some newer cells are characterized by a charge rate of C/3. These cells are capable of withstanding the higher charge rate indefinitely. The time required for fully charging such cells has thus been reduced to approximately 4 hours. However, even this amount of time may be too long for some applications.\nResearch by the battery manufacturers has also determined that properly designed cells may be charged at a rate of C/1, so that a cell may be recharged in approximately 1.2 hours, for the popular sub C (Cs) cell size. However, this approach can only be used if the high charge rate is terminated before the cells enter a destructive overcharge condition. For such cells, a maintenance, or \"trickle\" charge rate of C/10 is provided after the C/1 charge rate is complete. The trickle charge rate effectively \"tops off\" the battery charge and maintains the cell at full capacity until used.\nIt is thus necessary accurately to determine the particular point at which the permissible charge rate drops from the fast, C/1, rate to the trickle, C/10, rate in order to use the newly developed cells. Moreover, it is necessary to develop a control device which can accurately detect the changeover point and vary the charge rate accordingly.\nResearch into various cell characteristics which can be used for detecting the proper termination point for the C/1 charge rate has centered on voltage profiles, temperature changes, and internal pressure changes responsive to the charged state of the cell. Some prior art attempts have been directed to the use of internal cell pressure as the charging criterion. However, a special cell construction is required for sensing the internal pressure of the cell, involving access to the interior of the cell. The pressure sensing approach has thus not been widespread and is generally considered expensive.\nOther cell characteristics which have been considered as the criteria for determining the permissible charging state of the cell have included voltage and temperature.\nSensing the voltage alone, however, has generally not been found useful, since the voltage change from a discharged state to a fully charged state of the cell is small and is hard to detect accurately. More specifically, the change in voltage is typically of the same order of magnitude as the variation in voltage which may be found between cells of a battery. Such a variation, when within established tolerance levels, is small relative to the total cell voltage. Prior charging circuit designs have thus combined voltage sensing with temperature sensing, usually by placing a thermistor into intimate contact with the battery pack. However, the prior circuits, while generally effective, were complicated and expensive.\nMore recent improvements in cell design have made it possible to sense only the cell temperature as the criterion for terminating a C/1 charge rate. It is considered acceptable in the battery art to protect the cells from temperatures in excess of 45.degree. C. Thus, in known circuits thermostatically controlled switches are provided in intimate contact with the batteries. The thermostat is designed to open the associated switch at a temperature of 45.degree. C.\nA simple approach is used in one temperature sensitive arrangement of the prior art. Therein, the thermostatically controlled switch itself is used to break the fast charge current directly. A limiting resistor is provided in parallel with, and in close proximity to, the thermostat for supplying the C/10 maintenance charge current. In such a circuit, it is necessary to prevent further rapid charging of the battery cell once the trickle charge state has been entered.\nMore particularly, once the thermostatic switch has opened the rapid charge circuit the battery cell will cool, tending to reclose the thermostat and to reinitiate the process. Thus, it is necessary to latch the thermostatically controlled switch to an open condition once the fully charged state has been reached. The above described prior art approach utilizes the maintenance charge current to heat the thermostat, thereby to keep the thermostatic switch open for so long as the maintenance charge is continued. More specifically, in this approach the maintenance charge current is used to heat the limiting resistor for the trickle charging current. The close proximity between the limiting resistor and the thermostat provides a heat transfer therebetween, causing a temperature increase at the thermostat and opening the switch controlled thereby.\nAlthough the above concept is low in cost, such an approach requires continued heating of the thermostat by the C/10 limiting resistor. Since the thermostat is in intimate contact with the battery cells, however, the above described approach provides continued heating of the battery cells during the maintenance charging state. Such heating can shorten battery life. Moreover, the above described circuit leads to reduced reliability of the thermostatic switch, since the thermostat itself is required to break the large rapid-charge current at each termination of the rapid charge condition.\nIn another example of this approach, wherein the thermostat is required to break large currents in the rapid-charge mode, a gate of an SCR is biased by a capacitor and the thermostatic switch is in series with the SCR. It is thus necessary to control precisely the voltage on the capacitor in order to assure proper biasing of the SCR gate. Reliability of this approach suffers still further because of possible variations in capacitor parameters, and because of the difficulty of providing a more precise point at which to turn on the SCR.\nA more reliable concept has been to use the thermostat as a sensor only. In this approach, the thermostat is used to control associated electronics which, in turn, regulate the current. As with the previous approach, however, it is necessary to avoid overcharging the battery by a condition in which the fast charge rate is restarted once the batteries cool in the maintenance charge condition and the thermostat closes.\nThe major advantage of such an arrangement is that it is not necessary to heat the thermostat (hence the batteries) to latch the charger out of the fast charge mode while continuing a maintenance charge, since the charge rate is electronically controlled. Moreover, the thermostat is only required to switch a very low level sensing current rather than the full fast charge current.\nLow cost circuits utilizing the above approach are sensitive to one or more variables, however, such as battery or electronic component tolerance or battery impedance, which affects the reliability of latching the circuit out of the fast charge mode. In one such circuit the collector-emitter junction of a transistor is used to clamp across a gate-cathode junction of a power SCR. Such an arrangement does not necessarily keep the SCR off and is subject to variations in junction voltages of the transistor. Under particular circumstances it is thus possible that the SCR, supplying a high rapid charge current, may not be fully turned off and may overcharge the battery. Other circuits, using integrated circuits, controlled tolerance electronic devices, or other special techniques have been used to increase the reliability of the above described approach. However, such circuits are more expensive and thus are less desirable.\nThere is thus a need in the prior art for an inexpensive circuit, providing reliable recharging of battery cells and including reliable, low cost, control circuitry to avoid overheating and overcharging the battery.\nIt is accordingly an object of the invention to overcome the difficulties of the prior art and to provide a battery charging apparatus for rapid charging and maintenance charging of a battery.\nIt is a more specific object to provide a low cost battery charging apparatus wherein a thermostatic switch detects an appropriate transition point for terminating rapid charging and for initiating maintenance, or trickle charging of a battery.\nIt is still another object to provide a low cost battery charging apparatus utilizing a thermostatic switch to switch a low level sensing current rather than the full charging current.\nYet another object of the invention is the provision of a dual mode battery charging apparatus wherein a thermostatic switch senses an increased temperature to transfer charging from a rapid mode to a trickle mode and wherein a voltage providing circuit is used to latch the apparatus to the trickle mode when the thermostatic switch returns to a low temperature status.\nStill a more specific object of the invention is the provision of a voltage storage device for triggering a gate controlled SCR for rapidly charging a battery, including a circuit arrangement for changing the voltage level provided to the storage device in response to a temperature condition of the battery.\nYet a more particular object is an arrangement wherein an inverting structure is interposed between a voltage storage device and a gate controlled device triggered thereby, so that when a thermostatically controlled switch responds to a high temperature, fully charged, condition of a battery and increases the voltage of the storage device the reduced voltage of the inverting structure maintains the gate controlled device inactive even after the thermostatically controlled switch returns to a low temperature condition.\nIt is still a further object of the invention to provide a triggering device for a gate controlled device in a battery charging apparatus wherein a separate switch is required to be activated, in addition to activation of a thermostatically controlled switch, in order to cause a rapid charging operation and wherein reactivation of the thermostatically controlled switch, alone, will not reinitiate the rapid charging operation."}
{"text": "This invention relates to packaging and, more particularly, to a strap connector and strap system for package strapping and the like.\nHeretofore, myriad arrangements have been utilized and proposed for wrapping, tying, and strapping of packages, boxes, crates, and cartons. For example, it is common practice to package containers by using steel strapping, the straps being retained by crimped metal bands and by using fiberglass strapping of comparable character. Such strapping has constituted a great improvement in packaging as compared with the ubiquitous use of wire, cord, gummed, string or filament tapes, and other age old packaging expedients.\nHowever, metal and fiberglass strapping systems require use of special tools and techniques for proper application and tensioning which are not always available to the small commercial user or individual who wishes to take advantage of the security and tensile strength of such strap materials for packaging, tying down loads, fastening of loads, etc.\nAccordingly, it is an object of the invention to provide a strapping and strap connecting arrangement, including an improved strap and strap connector system.\nIt is another object of the invention to provide such a system which facilitates simple, facile, manual strapping, fastening, and securement of packages and other objects.\nIt is a still further object of the invention to provide such an arrangement which does not require the use of special tools or techniques for applying to a package or object secure, strong, tensile strapping.\nAnother object of the invention is the provision of such a connector which quickly and easily engages the ends of the strap for reliably and conveniently engaging the ends of such strap.\nYet another object of the invention is the provision of such a strap and strap connector which is readily constructed entirely of molded synthetic resin material by economic, mass production techniques.\nVarious other objects and features will be in part apparent and in part pointed out hereinbelow."}
{"text": "The human body posses the ability to resist most types of organisms, microorganisms or toxins that can cause damage to tissues and organs. This capacity is referred to as immunity. Much of the immunity is caused by an immunity system that forms antibodies and/or sensitized lymphocytes that attack and destroy the organisms, microorganisms or toxins.\nThis type of immunity is commonly referred to as acquired immunity.\nHowever, an additional portion of the immunity results from general processes rather than from processes directed at specific disease organisms. This is called innate immunity. Examples of innate immunity include Phagocytosis of bacteria and other invaders by white blood cells and reticuloendothelial cells; destruction of organisms swallowed into the stomach by the acid secretions of the stomach and by the digestive enzymes; resistance of the skin to invasion by organisms; and presence in the blood of special chemical compounds that attach to foreign organisms or toxins and destroy them.\nIn addition to its innate immunity, the human body also has the ability to develop extremely powerful specific immunity against individual invading agents such as lethal bacteria, viruses, toxins, and foreign tissues from other animals. This is called acquired immunity or adaptive immunity.\nThis system of acquired immunity is important as a Protection against invading organisms to which the body does not have innate immunity. The body does not block the invasion upon first exposure to the invader. However, within a few days to a few weeks after exposure, the special immune system develops extremely powerful resistance to the invader. Furthermore, the resistance is highly specific for that Particular invader and not for others.\nAcquired immunity can often bestow significant and long term protection. This is the reason the process known as vaccination is so important in the protection against disease.\nTwo basic, but related types of acquired immunity occur in the body. In one of these the body develops circulating antibodies, which are globulin molecules that are capable of attacking the invading agent. This type of immunity is called humoral immunity. The second type of immunity is achieved through the formation of large numbers of highly specialized lymphocytes that are specifically sensitized against the foreign agent. These sensitized lymphocytes have the special capability to attach to the foreign agent and to destroy it. This type of immunity is called cellular or cell-mediated immunity or, sometimes, lymphocytic immunity.\nBoth the antibodies and the sensitized lymphocytes are formed in the lymphoid tissue of the body. The presence of antigens initiate the immune process.\nAcquired immunity is the product of the body's lymphoid tissue. The lymphoid tissue is located mostly in the lymph nodes, but it is also found in special lymphoid tissue such as that of the spleen, in submucosal areas of the gastrointestinal tract, and, to a slight extent, in the bone marrow. The lymphoid tissue is distributed advantageously in the body to intercept the invading organisms or toxins before they can spread too widely.\nThough most of the lymphocytes in normal lymphoid tissue look alike when studied under the microscope, these cells are distinctly divided into two separate populations. One of the populations is responsible for forming the sensitized lymphocytes that provide cellular immunity and the other for forming the antibodies that provide humoral immunity.\nBoth of these types of lymphocytes are derived originally in the embryo from lymphocytic stem cells in the bone marrow. The descendants of the stem cells eventually migrate to the lymphoid tissue. Before doing so, however, those lymphocytes that are eventually destined to form sensitized lymphocytes first migrate to and are preprocessed in the thymus gland, for which reason they are called \"T\" lymphocytes. These are responsible for cellular immunity.\nThe other population of lymphocytes--those that are destined to form antibodies--is processed in some unknown area of the body, possible the liver and spleen. However, this population of cells was first discovered in birds in which the preprocessing occurs in the bursa of Fabricius, a structure not found in mammals. For this reason this population of lymphocytes is called the \"B\" lymphocytes, and they are primarily responsible for humoral immunity.\nTo further emphasize that the two populations of lymphocytes are separate, it is known that they also tend to localize in separate parts of the lymphoid tissue. For instance, in the lymph nodes, the lymphocytes of the \"B\" system are mainly located in the cortical and germinal areas, whereas the \"T\" cells are located in the paracortical areas.\nAfter formation of processed lymphocytes in the thymus, these first circulate freely in the blood and gradually filter into the tissues. Then they enter the lymph and are carried to the lymphoid tissue. The lymphoid tissue contains reticulum cells that form a fine reticulum meshwork. This filters the lymphocytes from the lymph, thereby entrapping them in the lymphoid tissue. Thus, the lymphocytes do not originate Primordially in the lymphoid tissue, but instead are transported to this tissue by way of the preprocessing areas of the thymus and probably the fetal liver.\nWhen a lymphocyte in the lymphoid tissue is stimulated to form either sensitized lymphocytes or antibodies, it always forms a sensitized lymphocyte or an antibody having specificity for a particular antigen. Because it is known that the lymphocytes of the lymphoid tissue can form literally hundreds or thousands of different types of sensitized lymphocytes and antibodies specific for different antigens, it is also almost certain that literally hundreds or thousands of different types of precursor lymphocytes pre-exist in the lymph nodes for formation of the many specific types of lymphocytes or antibodies.\nThe lymphocytes of each specific type in the lymphoid tissue--those that form one specific type of sensitized lymphocyte or one specific type of antibody--are called a \"clone of lymphocytes.\"\nEach clone of a lymphocyte is responsive to only a single type of antigen (or to a group of antigens that have almost exactly the same stereo-chemical characteristics). When excited by the clone's specific antigen, all the cells of the clone proliferate madly, forming tremendous numbers of progeny, and these in turn lead to the formation of large quantities of antibodies if the clone is \"B\" lymphocytes, or to the formation of numerous sensitized lymphocytes if the clone is \"T\" lymphocytes.\nOn exposure to proper antigens, sensitized lymphocytes are released from lymphoid tissue in ways that parallel antibody release except that instead of releasing antibodies, whole sensitized lymphocytes are formed and released from the lymphoid tissue into the lymph.\nAn important difference between cellular immunity and humoral immunity is its persistence. Humoral antibodies rarely persist more than a few months, or at most a few years. On the other hand, sensitized lymphocytes probably have an indefinite life span and seem to persist until they eventually come in contact with their specific antigen. There is reason to believe that such sensitized lymphocytes might persist as long as ten years in some instances.\nAlthough the humoral antibody mechanism for immunity is especially efficacious against more acute bacterial diseases, the cellular immunity system is activated much more potently by the more slowly developing bacterial diseases such as tuberculosis, brucellosis, and so forth. Also, this system is active against neoplastic cells, cells of transplanted organs, and fungal organisms, all of which are far larger than bacteria. And, finally, the system is very active against some viruses.\nTherefore, cellular immunity is especially important in protecting the body against some viral diseases, in destroying many early cancerous cells before they begin to grow, and, unfortunately, in causing rejection of tissues transplanted from one person to another.\nThe sensitized lymphocyte, on coming in contact with its specific antigen, combines with the antigen. This combination in turn leads to a sequence of reactions whereby the sensitized lymphocytes destroy the invader. As is also true of the humoral immunity system, the sensitized lymphocyte destroys the invader either directly or indirectly.\nSensitized lymphocytes can become bound with antigens in the membrane of an invading cell such as a cancer cell, a heart transplant cell, or a parasitic cell of another type. The immediate effect of this attachment is swelling of the sensitized lymphocyte and cause immediate release of cytotoxic substances from the lymphocyte to attack the invading cell.\nWhen sensitized lymphocytes combine with their specific antigens, they can release a number of different substances into the surrounding tissues that lead to a sequence of reactions. These reactions in turn are much more potent than the original attack on the invader. These include release of transfer factors which \"recruit\" additional lymphocytes having the same capability for causing the same cellular immunity reaction as the originally sensitized lymphocytes and includes the attraction and activation of macrophages to enhance their phagocytic and bacteriacidal functions.\nIt is by a combination of a weak direct effect of the sensitized lymphocytes on the antigen invader and much more powerful indirect reactions that the cellular immunity system destroys the invader.\nObviously, if a person should become immune to his own tissues, the process of acquired immunity would destroy his own body. Fortunately, the immune mechanism normally \"recognizes\" a person's own tissues as being completely distinctive from those of invaders, and his immunity system forms neither antibodies nor sensitized lymphocytes against his own cells and tissues. This phenomenon is known as tolerance to the body's own tissues.\nSeveral researchers have described direct cell-mediated cytotoxicity against tumors. These include, for example, descriptions of T-cell killing, (Vanky F, Klein E., \"Human T-cell cultures with selective antitomor reactivity,\" Cancer Immunol Immunother 1982;14:73-7; Vose BM, Bonnard GD, \"Specific cytotoxicity against autologous tumor and proliferative responses to human lymphocytes grown in interluekin 2.\" Int. J. Cancer 1982;29L33-9; Wunderlich J., \"Short-term .sup.51 Cr-release tests for direct cell-mediated cytotoxicity: methods, clinical uses and interpretations.\" In: Bach FH, Good RA, eds. Clinical immunobiology. New York: Academic Press, 1976:133-48); killing by macrophages, (Hibbs JB. \"The macrophage as a tumoricidal effector cell: a review of in vivo and in vitro studies on the mechanism of the macrophage non-specific cytotoxic reaction.\" In: Fink MA, ed. The Macrophage In Neoplasia. New York: Academic Press, 1976:83-112); natural cytotoxicity, (Stutman O, Lattime EC. \"Natural cell-mediated cytotoxicity against tumors in mice: an heterogeneous system, \" Transplant Proc. 1981;13:752-5); polyinosinic acid-induced cytotoxicity, (Dorfman N, Winkler D, Burton RC, Kassayda N, Sabia P, Wunderlich J., \" Broadly reactive murine cytotoxic cell induced in vitro under syngeneic conditions,\" J. Immunol. 1982;129-1762-9); lectin-induced cytotoxicity, (Mazumder A. Grimm EA, Zhang HZ, Rosenberg SA, \"Lysis of fresh human solid tumors by autologous lymphocytes activated in vitro with lectins. Cancer Res. 1982;42:913-8); and interleukin induced cytotoxicity, (Grimm EA, Mazumder A, Zhang HZ, Rosenberg SA, Lymphokine-activated killer cell phenomenon,\" J. Exp. Mec. 1982;55: 1823-41). Lectin-induced cytotoxicity and interleukin induced cytotoxicity represent polyclonal activation of T-cells.\nIt is widely known that spontaneously arising tumor cells are usually not sufficiently antigenic to induce an immune response. (Eggers et. al. Cancer Immunol. Immunother. (1982)).\nAdjuvant peptides have been used to stimulate the immunogencity of tumors. Adjuvants in general, are substances which non-specifically enhance the immune response to an antigen. It is known that in certain instances, administration of an antigen together with an adjuvant can completely change the mode of response, i.e., it is possible to break self tolerance to a large number of self or syngeneic antigens by injecting them into the host animal in an appropriate adjuvant. The most frequently utilized adjuvants are water-in-oil emulsions with the antigen in the aqueous phase, for example, Freund's Adjuvant. The adjuvant properties can be further enhanced by the addition of microbial antigen to the mixture such as in Freund's complete Adjuvant which contains heat-killed Mycobacterium tuberculosis.\nAntibody responses to antigens in adjuvants have proven to be greater, more prolonged and consist of different classes to the response obtained without an adjuvant. Adjuvants are hypothesized to work in a variety of ways. Initially, an antigen in an emulsion is resistant to dispersal and it therefore acts as a reservoir for antigen stimulation for-a period of long duration. Moreover, microbial products in general activate macrophages which lead to the Production of antigen non-specific factors which enhance the response. (Ivan M. Roitt et. al. IMMUNOLOGY Gower Medical Publishing, New York 1985).\nIn order to assist the body's own immune system in combating neoplastic diseases, such as all solid and lymphoid tumors, in vitro and in vivo immunization against tumor cell antigens has been performed.\nIn contrast to the non T-cell and polyclonal T-cell systems, it has been demonstrated that the use of an adjuvant peptide (AP) such as N-acetyl muramyl-L-alanyl-D-isoglutamine, covalently bound to the surface of poorly antigenic tumor cells such as murine methylcholanthrene-induced sarcoma cells (MC-1 cells) increases tumor immunogencity by permitting the induction of direct cell-mediated cytotoxicity against tumor associated antigens. This immunity is thus antigen driven and T-cell mediated.\nThe cell mediated cytotoxicity against tumor cells antigens has been shown to also be directed against cross reacting \"targets\" including many non H-2 (H=histo compatibility) matched solid tumors and some normal cells. This technique has been shown to work in vitro and in vivo, and a therapeutic effect has been demonstrated in vivo with small tumors. (Eggers, et. al. \"Use of Covalently Bound Adjuvant Peptide to Increase Tumor Immunogenicity\", Cancer Immunol Immunother; 12:167-172, 1982; Eggers, et. al. \"T-cell nature of Adjuvant Peptide-Induced Antitumor cell-mediated cytotoxicity\", J. Biological Response Modifiers, 3:387-390, 1984, Eggers et. al. \"In vivo immunization against autologous glioblastoma-associated antigens, Cancer Immunol. Immunother. 19:43-45 (1985) the disclosures of which are incorporated by reference herein).\nIt has been demonstrated that experimental animals as well as humans can elicit an immune response against their own neoplastic cells. It has been, therefore, somewhat of an enigma why, when humans are capable of such as immune response, in most cases, the tumor continues to grow and or metastasize and finally kills the host instead of being destroyed by an immune reaction of the host. An explanation for this failure has been attempted by David Ilfeld et. al. \"In vivo cytotoxicity and in vivo tumor enhancement induced by mouse spleen cells autosensitized in vitro,\" Int. J. Cancer 12:213-22, 1973.\nIlfeld et. al. hypothesized that circulating antibodies equipped with immunological specificity against tumor antigens accelerate tumor growth by actually interfering with the directed against tumor-specific antigens may favor the acceptance of an antigenic homograft whereas a strong response could lead to tumor rejection. Ilfeld et. al. studied the immuno-reactivity of C57B2 mouse spleen cells previously sensitized in vivo against syngeneic (self) fibroblasts, which were irradiated prior to sensitization. The sensitized spleen cells were assayed in vitro to test their cytotoxicity against 3LL tumor cells. Further, the sensitized spleen cells were mixed with the tumor cells and injected into syngeneic recipient mice in order to assay, in vivo, the influence on tumor growth. The in vitro assay showed that the sensitized lymphoid cells were cytotoxic against fibroblasts as well as against the tumor cells. However, in vivo experiments established that lymphoid cells sensitized against syngeneic fibroblasts promoted the growth of the 3LL tumor.\nA problem with using tumor cells to prepare a vaccine is that most human tumors cannot be grown in tissue culture. For example, most common human tumors, such as colon, breast, non-oat cell of the lung and prostate are difficult to grow in sufficient quantities to prepare enough for a single injection.\nIt is an objective of the present invention to provide a vaccine composition comprising non-malignant cells, preferably syngeneic, non-malignant cells coupled, to an adjuvant for administration to a vertebrate to induce a cell mediated cytotoxic response which cross reacts with tumor cells so that an anti-tumor immunity is induced in the vertebrate.\nIt is a second objective of the present invention to provide for a method for inducing anti-tumor immunity by administering non-malignant cells, preferably syngeneic, non-malignant cells coupled to an adjuvant as antigens or immunogens to induce in vivo cytolytic activity, i.e., induce an autoimmune response against the non-malignant cells that cross-reacts against neoplastic or tumor cells and induces in vivo protection against tumor cells.\nIt is a third objective of the present invention to provide a method for the treatment of neoplastic disease in a vertebrate comprising administering a therapeutically effective amount of non-malignant cells coupled with an adjuvant."}
{"text": "1. Field of the Invention\nThe present invention relates to an implantable heart monitoring device, with which it is possible to monitor the heart condition. The invention also concerns a corresponding method.\n2. Description of the Prior Art\nSeveral different devices for monitoring the performance of a heart are known. Often these devices are also able to deliver stimulation pulses to the heart. The devices are often able to sense the electrical activity in the heart. It is also known to determine an impedance value measured between different electrodes positioned in or at the heart. It is also known to sense other physiological parameters, such as pressure, oxygen level etc.\nUS 2001/0012953 A1 describes bi-ventricular pacing. An impedance may be measured between electrodes on the right and the left sides of the heart. The variation of the impedance with time is detected. The detected impedance variation may be used in order to synchronise the contraction of the ventricles.\nUS 2001/0021864 A1 describes different manners of using the proximal and distal electrodes of different leads in order to inject a current and to measure an impedance. The measured impedance value may be used in order to maximise the cardiac flow.\nUS 2007/0049835 A1 relates to an implantable cardioverter-defibrillator or pacemaker whose standard circuitry is used to trend a physiological cardiac parameter using intra-cardiac impedance measurements.\nUS 2007/0100249 A1 describes an implantable medical apparatus for detecting diastolic heart failure, DHF. The apparatus includes circuitry for determining, as the DHF parameter, the time duration of a predetermined phase of diastole.\nUS 2007/0055170 A1 describes a device for detecting the state of a heart on the basis of intracardial impedance measurement. The device has an impedance measuring unit as well as an analysis unit, which is connected to the impedance measuring unit and is implemented to derive a cardiac function parameter from a time curve of the impedance ascertained using the impedance measuring unit. The analysis unit derives a cardiac function parameter characterizing the behaviour of a heart during the diastole.\nU.S. Pat. No. 6,314,323 describes a heart stimulator in which the cardiac output is determined by measuring the systolic pressure.\nThe article “Hemodynamic Effects of Tachycardia in Patients with Relaxation Abnormality: Abnormal Stroke Volume Response as Overlooked Mechanism of Dyspnea Associated with Tachycardia in Diastolic Heart Failure” by Dae-Won Sohn et al., Journal of the American Society of Echocardiography, February 2007, pp. 171-176, describes a comparative study of two groups of individuals: healthy individuals and individuals with stable relaxation abnormality. The article describes how left ventricular pressure and stroke volume varies for the two groups when the heart is paced with 80 beats per minute and 120 beats per minute."}
{"text": "Vehicle mounted cable reel handling apparatus adapted to be carried on a trailer vehicle are shown in U.S. Pat. Nos. 3,091,413 and 3,063,584. Reel handling apparatus has also been provided for pickup trucks as disclosed in U.S. Pat. Nos. 3,165,214; 3,184,082; 3,036,790 and 3,325,118. In the above patents, the apparatus generally includes vertically movable lift arms pivotally connected to the vehicle for engaging and transferring a single ground supported reel onto the vehicle for transport. The lift arms are engageable with the reel at all times even when the reel is in the transport position.\nIn some of the reel handling structures for trucks to handle a pair of reels for transport, the operating cylinders for the reel lift arms are arranged on the truck bed so as to appreciably limit the space for reel storage as appears in U.S. Pat. Nos. 2,876,916 and 3,902,612. The Anderson U.S. Pat. No. 3,625,380 and McVaugh U.S. Pat. No. 3,820,673 use front and rear pairs of lift arms, with a first reel lifted from the ground by rear arms being transferred to the front lift arms and carried thereon to a farward transport position. The second reel when lifted from the ground remains on the rear lift arms for transport in a position adjacent to and rearwardly of the first reel.\nAlthough transfer of a reel from the rear lift arms to the front lift arms was generally satisfactory, the double lift arm arrangement was relatively expensive and difficult to accommodate within the limited space requirements on the truck bed, especially as restricted with the growing demand for larger side mounted tool carrying compartment units. With the compartment units extended from the truck cab to positions over and behind the truck rear wheel and axle assembly space requirements for transporting the reels become more critical.\nThe Hall U.S. Pat. No. 3,902,612 partially solves this problem by using transversely spaced tiltable beams extended longitudinally of the truck for receiving a reel from a pair of rear lift arms. On a controlled downward and forward tilting movement of the beams the transferred reel is rolled by gravity action to a forward transport position. However, by virtue of the lift arms being actuated by cylinders mounted on the truck bed, the transverse distance between the beams is appreciably reduced. As a result the reel has a spindle of reduced length, when lifted from the ground, which is then replaced by a longer spindle before the reel can be supported on the beams. Hall, therefore, has no provision for side compartment units and requires a manual changing of reel spindles, and a manual actuation of the tiltable beams to roll a reel to a front transport position. These disadvantages of the Hall apparatus are eliminated by the apparatus of this invention."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a magnetron mounting within a microwave oven.\n2. Description of the Prior Art\nHistorically, magnetron mounting in such devices as microwave ovens has required four attachment studs or bolts with their cooperating nuts. In some cases the studs were on the magnetron while in other cases they were on an oven wall or on brackets. Such variation requires the stocking of several different configurations to meet all design and repair requirements. In addition, labor requirements are high due to the time consuming task of individually joining and tightening each fastener combination. During both manufacture and repair, space constraints complicate the mounting of the magnetron.\nSpace constraints are imposed by the common desire to make the unit as compact as possible. Efficiency and safety require that any oven chamber, such as the cooking cavity and waveguide, if any, be \"sealed\" and have as few discontinuities as possible. For this reason, mounting of a magnetron within a microwave oven has often employed a bracket which extends beyond the chamber in question and to, or into, an adjacent closely packed compartment."}
{"text": "1. Field\nEmbodiments relate to a battery module.\n2. Description of the Related Art\nRecently, a high output battery module using non-aqueous electrolyte of high energy density has been considered. The high output battery module may include a high capacity battery module formed by serially connecting a plurality of battery cells to be used for driving a device that requires high power, e.g., a motor of an electrical vehicle.\nTypically, the battery module may include a plurality of battery cells. The battery cell may provide energy to an external electronic device by an electro-chemical reaction. The plurality of battery cells may be fixed (by a housing) to be used as a power source. The battery cell may include a highly reactive material therein, e.g., lithium. Thus, stability may be important."}
{"text": "Generating high quality and trusted random numbers is an essential task in various cryptographic schemes and many other applications such as Monte Carlo simulations [1] and various algorithms utilizing random numbers. Algorithmically generated pseudo-random numbers are available at very high rates and can be easily implemented in software, but they are deterministic in nature and therefore are not suitable for cryptographic purposes, and may be problematic for some Monte-Carlo simulations as well as other applications, e.g. games.\nAs an alternative, hardware random number generators have been used [2, 3]. They measure noisy physical processes and convert the outcome into random numbers. Since it is impossible to predict the outcome of such measurements, these physically generated random numbers are more trusted compared to pseudo-random numbers. Quantum random number generators (QRNG) are a class of hardware random number generators whose source of randomness is the outcome of measurements of a quantum noise source.\nEarly implementations of QRNGs made use of the decay statistics of radioactive nuclei [4, 5]. A number of more recent implementations using quantum optical measurements have been reported. These include measuring photon number statistics [6-12], scattering events of single photons by a beam splitter [13] and amplified spontaneous emission of a fiber amplifier [14]. QRNGs based on measuring the intensity [15, 16] and phase noise [17-24] of different light sources have also been reported.\nQRNG implementations based on measuring the vacuum fluctuations of the electromagnetic field have also been reported in [25-28]. Such measurements are known for their high bandwidth. However, the reported QRNG implementations based on measuring vacuum fluctuations of the electromagnetic field have commercial implementation problems due to complexity of the alignment of optical elements and issues in manufacturability.\nEmbodiments of the present invention provide an alternative method and system for random number generation."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus, a method and a system for resistance-based evaluation of a muscular force.\n2. Description of the Background Art\nA system for resistance-based evaluation of a muscular force, exploiting a bi-articular link mechanism, such as a bi-articular arm apparatus, has been proposed in Japanese patent laid-open publication No. 2000-210272, for instance. In this known system for resistance-based evaluation of a muscular force, the muscular output of the antagonistic mono-articular muscles and the antagonistic bi-articular muscles of a test subject, or testee, whose muscular force is being evaluated, is measured with a pressure sensor. The test subject, or person, exerts the force with maximum isometric effort, at his or her four limbs, in a plurality of preset directions. Then, a hexagonal diagram showing characteristics of the output distribution is prepared to use the so prepared diagram to evaluate the force of function-based muscles of praxis.\nThere has also been proposed a model for bi-articular muscles that functions for a mammal inclusive of the human being to bend an arm, and researches are now underway with the use of the model to control the movements of the bi-articular link mechanism. See T. Fujikawa et al., “Functional Coordination Control of Pairs of Antagonistic Muscles”, Transactions of The Japan Society of Mechanical Engineers, Vol. 63, No. 607 (1997-3) pp. 769-776, Treatise No. 96-1040. In this research, it is stated to be desirable to use, as a driving source, a model of an actuator having an elastic element and a contractile element that exerts the force in the contracting direction.\nIn the above-described conventional solutions for resistance-based muscular force evaluation, the four limbs of a test subject are simulated to a link system, and a muscular force is evaluated based on a hexagonal diagram plotting characteristics of output distribution. Since the hexagonally-shaped output distribution, plotted on the diagram, differs with the joint angles of the link system, measurement cannot be conducted unless proportionally contracted. Further, it is difficult to arrive at proper comparison if the test subject changes his or her position each time his or her muscular force is measured."}
{"text": "This invention relates to a product for use on the sand shore of a beach, and more particularly concerns a product which serves both as a beach blanket and chair.\nThere is widespread interest in recreational or leisure time visits to sandy beaches. For more comfortable relaxation at the beach, a beach \"blanket\", usually a large sheet or blanket structure, is spread upon the sand so that the beachgoer can sit or lie upon the beach surface without actually being in contact with the sand.\nAnother item generally brought to the beach is a beach \"chair\", usually of foldable light-weight construction. The beach chair provides support for the user's back, and is therefore more comfortable than sitting directly upon the sand or the beach blanket, especially if the beachgoer intends to read or converse.\nAlthough the beachgoer will want to bring to the beach equipment, food, and other articles to enhance the enjoyment of the occasion, the carrying and transporting of such items to the beach can be burdensome and diminish the enjoyment of the occasion.\nThe use of multifunctional beach related items are therefore of interest to the beachgoer, in that they minimize the amount of time and effort required to prepare and pack an automobile or bicycle for the trip to the beach, and in carrying the items onto the beach. Many types of compact beach chairs have earlier been disclosed, and various types of beach blankets are known, including inflatable structures which doubly serve as a raft. Inflatable chairs and related devices which provide support for the user's back are disclosed in U.S. Pat. Nos. 2,612,645; 3,112,956; 3,408,107, and 4,189,181. However, such prior devices are unsuitable for use on a beach or are not properly designed to support the back on a sand surface.\nIt is accordingly an object of the present invention to provide a product useful as both a beach blanket and beach chair, capable of accommodating one or more persons.\nIt is a further object of this invention to provide a product as in the foregoing object, a portion of which can be controllably inflated to form a comfortable chair-like structure having back-supporting characteristics.\nIt is another object of the present invention to provide a product of the aforesaid nature of rugged, durable, light-weight construction and amenable to low cost manufacture.\nThese objects and other objects and advantages of the invention will be apparent from the following description."}
{"text": "Three-dimensional (3D) scanning systems are used to capture 3D models of real-world, physical objects by using one or more 3D scanners to capture views of the physical objects from multiple angles and synthesizing the 3D models from the multiple views.\nThree-dimensional scanning systems may be contrasted with panoramic photography or virtual reality (VR) photography that capture a series of two dimensional (2D) images from multiple viewpoints and combine the 2D images into a format that allows a user to control the panning of a virtual camera across a 2D scene that remains stationary at a fixed central point (e.g., “outward facing” VR photography) and formats that allow a user to rotate around an object to view 2D images of the object from different sides (e.g., “inward facing” VR photography).\nIn contrast, the 3D scanning systems may be used to capture three-dimensional models of objects and/or environments (e.g., geometry defined by collections of three dimensional points), thereby allowing, for example, the measurement of lengths within the 3D models, and the insertion of the captured 3D models into 3D virtual environments such as 3D modeling tools and 3D game engines."}
{"text": "Membrane technology is extensively applied to gas separation. For this application, high gas permeability and gas pair selectivity are the two most important criteria for choosing membrane materials. See, e.g., P. Bernardo, et al., Ind. Eng. Chem. Res., 48, 2009, 4638 and V. Abets, et al., Adv. Eng. Mater. 8, 2006, 328.\nCross-linked aromatic polyimide materials have received much attention for use in gas separation. Methods to cross-link aromatic polyimides include ultraviolet irradiation of benzophenone-containing polyimide chains, thermal treatment of polyimide chains containing acetylene end groups, and cross-linking polyimide chains by small molecules with multiple reactive groups. See H. Kita, et al., J. Membrane. Sci. 87, 1994, 139; Y. Xiao, et al., J. Membrane. Sci., 302, 2007, 254; and C. Staudt-Bickel et al., J. Membr. Sci., 155, 1999, 145. These methods tend to increase chain packing and inhibit intra-segmental and inter-segmental mobility among chains, resulting in improved gas pair selectivity but sacrificing gas permeability.\nThere is a need for improved polyimide membranes that can be used in gas purification with both high gas permeability and gas pair selectivity."}
{"text": "Data is frequently transmitted to and from users, such as consumers, over a network such as the Internet, cable, fiber, wireless, or satellite networks. A data modem is frequently used to modulate and demodulate data for transmission over the network. Data modems normally include one or more packet buffers to store incoming and outgoing packets, which may be transmitted using protocols such as UDP (User Datagram Protocol), which generally provides real-time non-guaranteed delivery of packets, and TCP/IP (Transmission Control Protocol/Internet Protocol), which generally guarantees delivery of packets. Upstream packets such as UDP packets containing voice data or TCP/IP packets containing uploaded pictures or videos are generally stored in the packet buffer until they can be accepted by the network.\nThe ability of a given network to accept packets from the packet buffer may depend on a service “tier” to which a given consumer or other user has access. Users who have access to higher levels of service may be provided with a higher bandwidth corresponding to an improved performance experience, whereas users who have access to lower levels of service may be provided with a lower bandwidth. These factors may affect the length of time that packets remain in the packet buffer. If the network is congested, packets may remain in the buffer for a longer period of time, leading to perceptible delays, especially for applications such as voice transmission.\nWhen a data modem is “provisioned” for a particular user or class of users, one of the parameters that may be set is the size of the packet buffer. By setting the packet buffer to a large size, some packets may stay in the buffer for a long period of time, creating perceptible delays. By setting the packet buffer to a small size, the buffer may fill up quickly, leading to an underutilization of the provisioned data rate. A default packet buffer size may be provided, which may be based on an assumption that the provisioned user will have a high tier of service corresponding to high bandwidth. Once the packet buffer size is set, it is generally not changed for the user. It would be desirable to allow more flexibility by allowing the buffer size to be adapted over time based on one or more factors."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital camera having a self-timer shooting function and an image display function.\n2. Description of the Related Art\nIn resent years, digital cameras have rapidly been becoming widespread, and various types of digital cameras are supplied to the market. Generally, a digital camera frequently has an image display device comprising a liquid crystal display (hereinafter, referred to as an LCD) or the like on its back. Because of the LCD, the photographer can confirm the image to be shot without viewing through the finder. Therefore, it is possible to shoot the subject at an extremely free angle. Moreover, a conventional digital camera is provided with a function to continue displaying the shot image on the image display device for a predetermined time every time one frame is shot, that is, a function to hold the shot image for a predetermined time. Because of this function, the photographer can check whether the shot image is desired or not without performing any complicated operations.\nWhether silver halide film cameras or digital cameras, cameras are frequently provided with a function to perform shooting by use of a self-timer (hereinafter, referred to as a self-timer function). According to the self-timer function which is used, for example, when the photographer himself or herself is the subject to be shot, shooting is on standby for a predetermined time after the depression of the release button, and shooting is performed after the predetermined time has elapsed. Because of this function, the photographer can shoot himself or herself by moving to the position of shooting within the predetermined time after depressing the release button.\nHowever, since the photographer is the subject when self-timer shooting is performed as mentioned above, it takes time for the photographer to return to the camera after shooting is performed. Therefore, even though the digital camera has the image display device and the function to display the shot image for a predetermined time, it frequently occurs that the display of the shot image ends before the photographer returns to the camera to check the display device on the back of the camera. That is, the photographer cannot check the shot image. Since recorded images can be read out from a recording medium and played back, the shot image can be checked by playing it back. However, complicated operations-such as switching to a reproduction mode and specification of the frame to be played back are necessarily performed every time, which is inconvenient.\nIn view of the above-mentioned problem, an object of the present invention is to provide a digital camera in which the photographer can easily check whether the shot image is desired or not without performing any complicated operations even when self-timer shooting is performed."}
{"text": "In recent years, extremely many kinds of organic photochromic dyes has been developed, and photochromic dyes capable of being obtained as commercially available products have been increasing. The application thereof to a lens for eye glasses also becomes popular along with the trend of plasticization on the marketplace, and a photochromic lens made of plastic, applying an organic photochromic dye comes onto the market for use in eye glasses.\nAs to methods for manufacturing a photochromic lens, there are disclosed (1) a method of coating a resin liquid containing a photochromic compound on a lens, heating the same to cause the photochromic compound to permeate the lens surface layer, after that, removing the coated resin film, and applying a curable film thereon (for example, see Patent Literature 1), and (2) a method of dissolving a photochromic compound into a lens coating liquid, and coating and curing the same on the lens surface (for example, see Patent Literature 2).\nHowever, in the method of (1), in order to obtain a sufficient photochromic density, it is necessary to cause the photochromic compound with a high concentration to permeate the lens surface and there is a problem in which a lens base material is limited to a material to be highly permeated. Therefore, in terms of heat resistance, mechanical strength and the like, a satisfactory level as a lens for eye glasses is not attained. In addition, in the method of (2), there is a limitation on the solubility of the photochromic compound into the coating liquid and the securement of a sufficient coloration density is difficult.\nFurthermore, in these methods, since a film is formed by coating the coating liquid on lens surfaces of variously curved shapes, a high-accuracy technology of homogenizing the film and a high-accuracy technology of controlling a film thickness, which correspond to these methods, are required, and thus the manufacturing cost becomes high.\nIn contrast, as manufacturing methods other than (1) and (2), (3) a method of dissolving previously a photochromic compound into a monomer mixed liquid for a lens, pouring it into a mold, and after that, polymerizing it to obtain a photochromic lens is disclosed (for example, see Patent Literature 3 and Patent Literature 4).\nMore particularly, in Patent Literature 3, a photochromic lens is disclosed, the lens having a sufficient light-controlling performance, and excellent surface hardness and abrasion resistance that are important as a lens. In addition, in Patent Literature 4, a photochromic lens is disclosed, the lens having a low yellow level before coloration and having wavelength in coloration made longer, to thereby allow a deep tone to be expressed.\nMethods of (1) and (2) require a special process for imparting light controllability such as a coating treatment after lens molding. In contrast to this, the method of (3) is preferable as a manufacturing method because the light-controlling performance is imparted simultaneously at the time of the lens molding and thereby manufacturing number of processes is lowered, and in addition, since the photochromic compound can easily be dispersed homogeneously in a base material, the method of (3) is extremely useful as a method for mass-producing a lens having a certain light-controlling performance irrespective of the lens shape and having a stable quality.\nFurthermore, as specific examples of the method of (3), in Patent Literature 5 and Patent Literature 6, it is described that the combination of a specific aromatic (meth)acrylic ester and aromatic vinyl makes it possible to obtain a good light-controlling performance."}
{"text": "Credit card sized integrated circuit integrated circuit cards are commercially available and find utility in numerous electronic systems. In personal computers they supplement or replace floppy disks by carrying software programs and data. Integrated circuit cards are connected to the computer's internal logic by insertion through an opening in the side of the computer to mate with a socket contained within the computer. Portable personal computers, notebook computers, and pocket diaries particularly find integrated circuit cards convenient because they avoid the need for the expense, power requirements, bulk, and weight of a disk drive; the integrated circuit card needs only an electrical connector and minimal structural, support.\nFacsimile and copy machines can use an integrated circuit card to store data related to usage control. Typewriters and printers can use an integrated circuit card to store desired memory fonts. Word processors can use an integrated circuit card to store text. Hand-held terminals can use integrated circuit cards to store inventory control information. Electronic cash registers can use integrated circuit cards to store price information. Controllable machinery can use integrated circuit cards to store automation control information. Programmable controllers can use integrated circuit cards to store process control data. Electronic game systems can use integrated circuit cards to store the specifics of games to be enacted on TV screens by the players.\nOther areas that can benefit from the use of integrated circuit cards include bulk data acquisition such as in music and photography, where the desired song or picture is stored in the memory devices of the card.\nPresently available integrated circuit cards typically include one or more plastic-encapsulated or other types of integrated circuits solder attached to connection stripes or spots on a printed wire board. A connector is solder attached to metal lines which are formed on the printed wire board and extend to one edge of the board. The connector connects external power, signal, and ground lines to circuitry in the card. An external shell or encasement surrounds the printed wire board/integrated circuit/connector assembly. The encasement typically has top and bottom frames formed from molded plastic and bonded together along an outer edge to define an interior chamber in which the printed wire board is located. The top and bottom frames have ribs which contact and securely hold the printed wire board and connector in place. Top and bottom metal plates are bonded to the top and bottom frames, respectively, and cover openings in the top and bottom frames.\nThe use of a plastic frame adds significantly to the complexity, tooling costs, and tooling turn around time in integrated circuit card construction. This is due to the fact that the manufacture of a plastic frame requires the use of a mold. Whenever the size of the printed wire board or the number, location, or size of the semiconductor devices on the printed wire board is changed, a new plastic frame that will accommodate the changes is required. The fabrication of the mold for the new frame is extremely time consuming and expensive. As a result, the time and expense incurred in producing integrated circuit cards incorporating even minor design changes are significant.\nThe use of a molded plastic frame is also a major source of problems during assembly of the integrated circuit card due to the fact that molded plastic frames are easily bent and often have surface irregularities introduced during the molding process which prevent all integrated circuit card parts from fitting together properly. In addition, processes for bonding the top and bottom plastic frames together, such as ultrasonic welding, are difficult and produce integrated circuit cards having inconsistent quality.\nAccordingly, a need exists for an integrated circuit integrated circuit card that can be produced quickly, is simple, inexpensive, and reliable, and does not require a plastic frame."}
{"text": "1. Field of the Invention\nThis invention relates to a suction cup device, more particularly to a suction cup device which defines a volume-variable space to produce an adjustable reduced pressure with a desired suction strength.\n2. Description of the Related Art\nFIG. 1 shows a conventional suction cup device 10 disclosed in U.S. Patent Publication No. US 2010/0116954 A1, which is adapted to be attached to a flat wall 1, and which includes a suction cup 13 with a rack 14 connected to the center thereof, a mount body 11 attached to the suction cup 13 by virtue of a pressing ring 12, and a lever 15 pivotably mounted on the mount body 11 and provided with a pinion 152 meshing with the rack 14. By turning the lever 15, a central part of the suction cup 13 can be pulled away from the flat wall 1 to produce a reduced pressure in an enclosed space between the flat wall 1 and the suction cup 13 for holding the suction cup device 10 against the flat wall 1. However, in the conventional suction cup device 10, the strength of suction force generated as a result of movement of the central part of the suction cup 13 is constant, and cannot be increased after a period of use, thereby resulting in undesired disengagement of the suction cup device 10 from the flat wall 1.\nAnother conventional suction cup device is disclosed in U.S. Pat. No. 7,021,593 B1, which includes a suction lock/release disposed on a lever and deep scoops disposed on the mount body to lock the lever at a fixed angle so as to adjust the suction strength of a suction cup. However, assembly of the suction lock/release to the lever is troublesome. Besides, manual operation of the suction lock/release is required to lock or unlock the lever, thereby rendering the adjustment complicated."}
{"text": "1. Field of the Invention\nThis invention relates to magnetic thin film heads (TFH) for recording and reading magnetic transitions on a moving magnetic medium.\n2. Background of the Invention\nMagnetic TFH transducers are known in the prior-art. See, e.g. U.S. Pat. Nos. 4,016,601; 4,190,872; 4,652,954; 4,791,719 for inductive devices and U.S. Pat. Nos. 4,190,871 and 4,315,291 for magnetoresistive (MR) devices.\nIn the operation of a typical inductive TFH device, a moving magnetic storage medium is placed near the exposed pole-tips of the TFH transducer. During the read operation, the changing magnetic flux of the moving storage medium induces changing magnetic flux upon the pole-tips and gap between them. The magnetic flux is carried through the pole-tips and back-portion core around spiralling conductor coil winding turns located between the core arms. The changing magnetic flux induces an electrical voltage across the conductor coil. The electrical voltage is representative of the magnetic pattern stored on the moving magnetic storage medium. During the write operation, an electrical current is caused to flow through the conductor coil. The current in the coil induces a magnetic field across the gap between the pole-tips. A fringe field extends into the nearby moving magnetic storage medium, inducing (or writing) a magnetic domain (in the medium) in the same direction. Impressing current pulses of alternating polarity across the coil causes the writing of magnetic domains of alternating polarity in the storage medium. Magneto-resistive (MR) TFH devices can only operate in the read mode. The electrical resistance of an MR element varies with its magnetization orientation. Magnetic flux from the moving magnetic storage medium induces changes in this orientation. As a result, the resistance of the MR element to a sensing electric current changes accordingly. The varying voltage signal is representative of the magnetic pattern stored on the magnetic medium.\nPrior-art magnetic recording inductive thin film heads include top and bottom magnetic core pole layers, usually of the alloy Ni--Fe (permalloy), connected through a via in the back-portion area, and separated by a thin gap layer between the pole-tips in the front of the device. The bottom pole-tip is usually designed to be wider than the top pole-tip in order to prevent \"wraparound\" due to misregistration or misalignment, as taught by R. E. Jones in U.S. Pat. No. 4,219,855. Alternatively, one or both pole-tips are trimmed by ion-milling or by reactive ion etching (RIE) to ensure similar width and proper alignment. Such a technique is disclosed, for example, by U. Cohen et al. in U.S. Pat. No. 5,141,623. As the track width decreases in order to increase the recording density, the write head pole-tips must be very narrow. P. K. Wang et al. describe elaborate schemes to obtain pole-tips for writing very narrow track width, in IEEE Transactions on Magnetics, Vol. 27, No. 6, pp. 4710-4712, November 1991.\nOne of the problems associated with the prior-art pole-tip designs is that during write operations, substantial noise is introduced along the track-edges (on the magnetic storage medium), which adds to the noise generated by the medium during read operations. During the write operations, significant portions of the intense magnetic flux lines, emanating from the corners and side-edges of the pole-tips, deviate from a direction parallel to the track's length. The non-parallel magnetic field magnetizes the medium in the wrong directions, giving rise to noise along the track-edges. This noise is usually characterized as \"track-edge fringing noise\" and is a major obstacle to increasing the track density. According to a paper by J. L. Su and K. Ju in IEEE Transactions on Magnetics, Vol. 25, No. 5, pp 3384-3386, September 1989, the track-edge noise extends about 2.5 .mu.m on each side of the written track. As track density increases, the track width decreases along with the strength of the read-back signal. If the track-edge fringing noise remains the same, then the signal to noise ratio (SNR) is directly proportional to the track width, and deteriorates rapidly as the latter decreases. The current state-of-the-art magnetic thin film media can support lineal density of about 40,000-60,000 flux changes per inch (FCI), corresponding to domain length of about 0.4-0.7 .mu.m. Yet, the track width is at least an order of magnitude larger, about 8-12 .mu.m. There is no apparent reason why the media could not support much narrower tracks, if not for the rapid deterioration of the SNR. By eliminating most of the track-edge fringing noise, the useful track width could be decreased to about 1.0 .mu.m, or less. This represents an increase of recording density by about an order of magnitude.\nIn addition to the medium's noise, there is also the head's noise. A significant portion of the head's noise is due to edge-closure domains in the pole-tips. This noise contribution becomes more dominant as the width of the pole-tips decreases. This problem was described, for example, by D. A. Herman in Paper No. 299, \"Laminated Soft Magnetic Materials\", The Electrochemical Society Conference, Hollywood, Fla., October 1989."}
{"text": "Circuit structures are known in which a conductive member that constitutes a circuit for conducting a comparatively large current is fixed to a substrate provided with a conductive pattern that constitutes a circuit for conducting a comparatively small current (for example, see JP 2003-164040A below).\nIn the circuit structure disclosed in JP 2003-164040A, a main portion of an electronic component (FET) is mounted on the conductive member, and at least one of the terminals of the main portion is connected to the conductive member while the other terminal or terminals are connected to the substrate (see FIG. 4 of JP 2003-164040A). Because there is a height difference (step difference), which corresponds to the thickness of the substrate, between a surface of the substrate and a surface of the conductive member, there is a need to perform processing such as bending of either of the terminals. Also, there may be cases where bending cannot be performed if the terminals are short, and connection becomes difficult.\nThe issue to be solved by the present invention is to provide a circuit structure to which an electronic component can be easily mounted (electrical connection of terminals)."}
{"text": "1. Field of the Invention\nThe invention relates to a process for measuring the flow vectors in gas currents containing optically detectable particles wherein a focusing means focuses at least two spatially separated beams in at least two focusing point in a measuring volume.\n2. Description of the Related Art\nProcesses for measuring flow vectors in gas currents have been known to employ light of a light source focused by a focusing means in the flow channel at two focusing points positioned in a close succession (U.S. Pat. No. 3,941,477). Particles contained in the gas current are illuminated in traversing the focusing points. Due to the stray radiation reflected by the particles, a start pulse is produced when the first focusing point is traversed, while a stop pulse is produced during the traversing of the second focusing point. From the time interval between the start pulse and the stop pulse, it is possible to determine the component of the particle speed vector in direction of the straight line traversing the focusing points. By moving the focusing device, said direction may be varied thus permitting the detection of flow vectors having different directions. However, by said method, it is only possible to measure the components of the flow vectors which extend in a normal plane relative to the optical axis of the focusing means. The component extending in the direction of the optical axis may not be determined. Said process is designated to 2d-process to refer to the two-dimensional vector measurement.\nA further development of said process is designated as 3d-process, by which the vector component extending in direction of the optical axis may be detected as well (British Pat. No. 2,109,548) and in which four laser beams are used two of which each form a beam pair. The beams of each pair being directed to two focusing points situated in the same normal plane relative to the optical axis. Due to the differences of the direction of incidence of the beams directed to a focusing point, the flow directions measured by means of the beam pairs are determined differently. From said difference of direction, one may draw a conclusion concerning the flow component in direction of the optical axis of the system. The expenditure and the laser capacity required by said known process are quite considerable."}
{"text": "Alternative transportation fuels have been represented as enablers to reduce toxic emissions in comparison to those generated by conventional fuels. At the same time, tighter emission standards and significant innovation in catalyst formulations and engine controls has led to dramatic improvements in the low emission performance and robustness of gasoline and diesel engine systems. This has reduced the environmental differential between optimized conventional and alternative fuel vehicle systems. However, many technical challenges remain to make the conventionally-fueled internal combustion engine a nearly zero emission system having the efficiency necessary to make the vehicle commercially viable.\nAlternative fuels cover a wide spectrum of potential environmental benefits, ranging from incremental toxic and carbon dioxide (CO2) emission improvements (reformulated gasoline, alcohols, etc.) to significant toxic and CO2 emission improvements (natural gas, etc.). Hydrogen has the potential to be a nearly emission free internal combustion engine fuel (including CO2 if it comes from a non-fossil source).\nThe automotive industry has made very significant progress in reducing automotive emissions. This has resulted in some added cost and complexity of engine management systems, yet those costs are offset by other advantages of computer controls: increased power density, fuel efficiency, drivability, reliability and real-time diagnostics.\nFuture initiatives to require zero emission vehicles appear to be taking us into a new regulatory paradigm where asymptotically smaller environmental benefits come at a very large incremental cost. Yet, even an “ultra low emission” certified vehicle can emit high emissions in limited extreme ambient and operating conditions or with failed or degraded components.\nOne approach to addressing the issue of emissions is the employment of fuel cells, particularly solid oxide fuel cells (SOFC), in an automobile. A fuel cell is an energy conversion device that generates electricity and heat by electrochemically combining a gaseous fuel, such as hydrogen, carbon monoxide, or a hydrocarbon, and an oxidant, such as air or oxygen, across an ion-conducting electrolyte. The fuel cell converts chemical energy into electrical energy. A fuel cell generally consists of two electrodes positioned on opposite sides of an electrolyte. The oxidant passes over the oxygen electrode (cathode) while the fuel passes over the fuel electrode (anode), generating electricity, water, and heat.\nThe fuel gas for the cell can be derived from conventional liquid fuels, such as gasoline, diesel fuel, methanol, or ethanol. The device, which converts the liquid fuel to a gaseous fuel suitable for use in a fuel cell, is known as a reformer.\nThe long term successful operation of a fuel cell depends primarily on maintaining structural and chemical stability of fuel cell components during steady state conditions, as well as transient operating conditions such as cold startups and emergency shut downs. The support systems are required to store and control the fuel, compress and control the oxidant and provide thermal energy management."}
{"text": "Pharmacokinetic and immune stimulating properties of proteins and synthetic drugs may be controlled by their conjugation to certain polymers. For example, polyethylene glycol (PEG) can be conjugated to proteins to achieve this effect (Fee and Van Alstine, Chemical Engineering Science, 61:924-934 (2006)). Such conjugation can take place if the relatively non-reactive hydroxyl groups present in PEG molecules are substituted by other, more reactive moieties (Jagur-Grudzinski, Reactive & Functional Polymers, 39:99-138 (1999)). A standard, linear PEG molecule is chemically a diol, which could suggest that the process of PEG derivatization and purification of products should be trivial. However, the polymeric nature of this diol, together with its amphiphilic properties can make these manipulations difficult. In some cases, the typical laboratory process for separation of difficult reaction mixtures, silica gel-based flash column chromatography, can fail for PEG with molecular weight higher than 1000. Neither crystallization nor precipitation appear adequate to achieve separation of PEG-containing materials, even if these methods can be used for efficient removal of other, contaminating substances with low molecular weight. Most reaction mixtures containing modified PEG molecules lack a reliable analytical method to control or to prove their composition. Polymers with functions that influence only minimally the hydrophobic properties of the polymer can be difficult to analyze by chromatography. The same applies for polymers with functions carrying only a minimal charge. This also applies for preparative chromatographic separation of charged polymers as described elsewhere for the separation of mono- and di-carboxyl modified PEG molecules (Drioli et al., Reactive & Functional Polymers, 48:119-128 (2001)).\nConfirmation of results of the synthesis based on NMR can be useless, as long as one is not sure about the purity of the product, and this is typically only obtained by chromatographic methods. This unusual conclusion comes from observations that an equimolar mixture of non-derivatized polymer and bis-derivatized polymer will produce an NMR pattern identical to the pure mono-derivatized polymer. Mass spectrometry can be complicated since most PEG exists not in the form of a single component, but is rather a Gaussian population of different polymer lengths, centered on its average molecular weight. Thus, even if all distinct components of the same type should have their mass increased by the same factor, the presence of unreacted and bis-modified material can obscure the picture of the analysis. The literature discusses this problem only sporadically, and often nothing is mentioned about analysis of the product or its purification. Many authors make the impression that the process that they describe is ongoing with quantitative yield, and thus the quality of the product does not need to be analyzed or questioned. This non-scientific approach can be frequently encountered in the chemistry of PEG. There are many examples in the literature presenting synthetic procedures with four to five consecutive steps without a single analysis of the product at any of these steps, without any attempts to purifying the product, and assuming 100 percent purity at the end of the process. It is, therefore, not strange that researchers after closer testing question these products and their purity (Ananda et al., Anal. Biochem., 374, 231-242 (2008)). A commonly accepted escape from the problem of selective modification is to work with a polymer that has one end blocked from the beginning by a stable chemical group, most often a methoxy group (mPEG). In theory, this blockage converts a PEG molecule to a monofunctional compound, and as such, it could be fully converted to the second derivatized form by increasing the amount of derivatizing reagent and/or time for reaction. Unfortunately, many of reactions commonly applied for derivatization of PEG are sluggish and only seldom go to completion. On the other hand, mPEG preparations contain significant percentages of PEG diol component. Moreover, the amount of this contamination increases with the length of mPEG, and this contamination can be hard to avoid. Consequently, derivatization will also result in formation of symmetrical, bis-derivatizated PEG, and its presence in the conjugating mixture results in formation of cross-linked products with unknown pharmacologic properties or a possible loss of protein activity. Therefore, pure, monofunctional polymers are usually preferred for protein modification, but one should be aware that purification of mPEG from its diol PEG contamination is practically impossible.\nNearly all of existing reactions, used today for derivatization of PEG, belong either to the alkylation-based or the acylation-based category. In the first case, the alkoxy anion, generated from PEG, is reacting with incoming electrophilic modifying reagent. Eventually, the activated PEG, subjected with a good leaving group, is itself an object of a nucleophilic attack. To this category belong processes resulting in thiolation, amination, azidation, and introduction of a carboxyl or an aldehyde group. Modified PEG's of this category will have their functional group connected directly to the PEG terminal carbon atom or these groups will be linked via an ether bond, a thioether bond, or a secondary amino group.\nThe second category, acylation, is based on a nucleophilic reaction of PEG's hydroxyl, (or another group present in a modified PEG—often an amino group), on an incoming acylating reagent. In many cases, this first acylation is followed by a second acylation that actually introduces the modification of interest to the PEG molecule. Functional groups incorporated by this method can be linked to the rest of PEG by an amido, a carbamido, urethane, thiourethane, or a simple ester group. These linking groups and the chemistry behind them belong to the very traditional methods of combining two chemical identities.\nPolyethylene glycols (PEG) coupled to phosphoramidites are used for direct coupling of PEG molecules to synthetic nucleic acids. One example is 4,4′-dimethoxy-trityl-polyethyleneglycol-[(2-cyanoethyl)-(N,N-diisopropyl)]-phosphoramidite. In these compounds, the phosphoramidite group is the part of the reactive functionality for linking the compound to a synthetic nucleic acid. It is designed to work in a completely water-free environment: In the presence of water, the phosphoramidite group can decompose instantaneously, making such PEG phosphoramidites inappropriate for conjugation to biological material in water-containing or aqueous solution. In particular, these PEG phosphoramidites can be inappropriate for conjugation to biological substances which are not soluble, stable or sufficiently reactive in non-aqueous media. Furthermore, already mildly acidic biological substances can decompose these PEG phosphoramidites. Finally, these PEG phosphoramidites contain a labile protecting group adjacent to the phosphorous atom which is specially designed to convert the intermediate phosphotriester to a phosphodiester. Phosphodiesters can be readily degraded enzymatically in vivo."}
{"text": "Many processes in biology, including transcription, translation, and metabolic or signal transduction pathways, are mediated by noN-covalently-associated multienzyme complexes1, 101. The formation of multiprotein or protein-nucleic acid complexes produce the most efficient chemical machinery. Much of modern biological research is concerned with identifying proteins involved in cellular processes, determining their functions and how, when, and where they interact with other proteins involved in specific pathways. Further, with rapid advances in genome sequencing projects there is a need to develop strategies to define “protein linkage maps”, detailed inventories of protein interactions that make up functional assemblies of proteins2,3. Despite the importance of understanding protein assembly in biological processes, there are few convenient methods for studying protein-protein interactions in vivo4,5. Approaches include the use of chemical crosslinking reagents and resonance energy transfer between dye-coupled proteins102, 103. A powerful and commonly used strategy, the yeast two-hybrid system, is used to identify novel protein-protein interactions and to examine the amino acid determinants of specific protein interactions4,6-8. The approach allows for rapid screening of a large number of clones, including cDNA libraries. Limitations of this technique include the fact that the interaction must occur in a specific context (the nucleus of S. cerevisiae), and generally cannot be used to distinguish induced versus constitutive interactions.\nRecently, a novel strategy for detecting protein-protein interactions has been demonstrated by Johnsson and Varshayskyl108 called the ubiquitin-based split protein sensor (USPS)9. The strategy is based on cleavage of proteins with N-terminal fusions to ubiquitin by cytosolic proteases (ubiquitinases) that recognize its tertiary structure. The strategy depends on the reassembly of the tertiary structure of the protein ubiquitin from complementary N- and C-terminal fragments and crucially, on the augmentation of this reassembly by oligomerization domains fused to these fragments. Reassembly is detected as specific proteolysis of the assembled product by cytosolic proteases (ubiquitinases). The authors demonstrated that a fusion of a reporter protein-ubiquitin C-terminal fragment could also be cleaved by ubiquitinases, but only if co-expressed with an N-terminal fragment of ubiquitin that was complementary to the C-terminal fragment. The reconstitution of observable ubiquitinase activity only occurred if the N- and C-terminal fragments were bound through GCN4 leucine zippers109,110. The authors suggested that this “split-gene” strategy could be used as an in vivo assay of protein-protein interactions and analysis of protein assembly kinetics in cells. Unfortunately, this strategy requires additional cellular factors (in this case ubiquitinases) and the detection method does not lend itself to high-throughput screening of cDNA libraries.\nRossi, F., C. A. Charlton, and H. M. Blau (1997) Proc. Nat. Acad. Sci. (USA) 94, 8405-8410) have reported an assay based on the classical complementation of α and ω fragments of β-galactosidase (β-gal) and induction of complementation by induced oligomerization of the proteins FKBP12 and the mammalian target of rapamycin by rapamycin in transfected C2C12 myoblast cell lines. Reconstitution of b-gal activity is detected using substrate fluorescein di-β-D-galactopyranoside using several fluorescence detection assays. While this assay bears some resemblance to the present invention, there are several significant distinguishing differences. First, this particular complementation approach has been used for over thirty years in a vast number of applications including the detection of protein-protein interactions. Krevolin, M. and D. Kates (1993) U.S. Pat. No. 5,362,625) teaches the use of this complementation to detect protein-protein interactions. Also achievement of β-gal complementation in mammalian cells has previously been reported (Moosmann, P. and S. Rusconi (1996) Nucl. Acids Res. 24, 1171-1172). The individual PCAs presented here are completely de novo designed interaction detection assays, not described in any way previously except for publications arising from applicants laboratory. Secondly, this application describes a general strategy to develop molecular interaction assays from a large number of enzyme or protein detectors, all de novo designed assays, whereas the β-gal assay is not novel, nor are any general strategies or advancements over previously well documented applications given.\nAs in the USPS, the yeast-two hybrid strategy requires additional cellular machinery for detection that exist only in specific cellular compartments. There is therefore a need for a detection system which uses the reconstitution of a specific enzyme activity from fragments as the assay itself, without the requirement for other proteins for the detection of the activity. Preferably, the assay would involve an oligomerization-assisted complementation of fragments of monomeric or multimeric enzymes that require no other proteins for the detection of their activity. Furthermore, if the structure of an enzyme were known it would be possible to design fragments of the enzyme to ensure that the reassembled fragments would be active and to introduce mutations to alter the stringency of detection of reassembly. However, knowledge of structure is not a prerequisite to the design of complementing fragments, as will be explained below. The flexibility allowed in the design of such an approach would make it applicable to situations where other detection systems may not be suitable.\nRecent advances in human genomics research has led to rapid progress in the identification of novel genes. In applications to biological and pharmaceutical research, there is now the pressing need to determine the functions of novel gene products; for example, for genes shown to be involved in disease phenotypes. It is in addressing questions of function where genomics-based pharmaceutical research becomes bogged down and there is now the need for advances in the development of simple and automatable functional assays. A first step in defining the function of a novel gene is to determine its interactions with other gene products in an appropriate context; that is, since proteins make specific interactions with other proteins or other biopolymers as part of functional assemblies, an appropriate way to examine the function of a novel gene is to determine its physical relationships with the products of other genes.\nScreening techniques for protein interactions, such as the yeast “two-hybrid” system, have transformed molecular biology, but can only be used to study specific types of constitutively interacting proteins or interactions of proteins with other molecules, in narrowly defined cellular and compartmental contexts and require a complex cellular machinery (transcription) to work. To rationally screen for protein interactions within the context of a specific problem requires more flexible approaches. Specifically, assays that meet criteria necessary not only to detecting molecular interactions, but also to validating these interactions as specific and biologically relevant.\nA list of assay characteristics that meet such criteria are as follows:\n1) Allow for the detection of protein-protein, protein-DNA/RNA or protein-drug interactions in vivo or in vitro.\n2) Allow for the detection of these interactions in appropriate contexts, such as within a specific organism, cell type, cellular compartment, or organelle.\n3) Allow for the detection of induced versus constitutive protein-protein interactions (such as by a cell growth or inhibitory factor).\n4) To be able to distinguish specific versus non-specific protein-protein interactions by controlling the sensitivity of the assay.\n5) Allow for the detection of the kinetics of protein assembly in cells.\n6) Allow for screening of cDNA, small organic molecule, or DNA or RNA libraries for molecular interactions."}
{"text": "1. Field of the Invention\nThe present invention relates generally to databases, and more particularly, to identifying and presenting fields from databases.\n2. Description of the Related Art\nEarly database software programs were a great relief for people who needed to organize and store vast amounts of data. Thus, databases allowed people to input and store information in a form that could be easily re-called and updated. As is well known in the art, database programs have gained tremendous acceptance and usage by families and students, as well as in business settings. By way of example, families have used database programs to keep medical records, keep a budget, maintain an inventory of assets, and compile wedding plan information. Similarly, students may use a database program to maintain student loan records, prioritize class events, and coordinate field trips.\nGenerally, database programs have gained increased acceptance by the computing public with the advent of more user-friendly database programs, which have made data entry more efficient. Although there have been great improvements in data entry, there is still a great need for database programs that reduce the hassles associated with viewing the data in various formats. By way of example, it may be desirable to view all data associated with a single record. As another example, it may be desirable to compare selected data associated with multiple records. Thus, each time a user wishes to view data in a different format, it is necessary to create a new “report template.” As a result, this typically requires repetitively entering frequently viewed fields from scratch each time the information is needed for a particular report template.\nIn operation, when a database designer wants to define fields for a particular database, a form layout window can be used. FIG. 1 shows a screen shot illustrating an exemplary form layout window used to define fields for a database. In this example, an optional header, “Golf Money Winners” 102, is created and a number of fields 104 are defined to create a body 103. As shown, a last name field 106, a first name field 108, an events field 110, and a total prize money field 112 are defined in the body 103 of the form layout window. Although the form layout may include a footer, a footer is not included in this example.\nOnce data has been entered and stored in the database, it is often desirable to view the data for a particular record. In the simplest case, the data may be viewed in the format in which it has been entered. Thus, the form layout window is suitable for obtaining (i.e., entering) data for a particular record as well as for displaying data associated with a single database record. FIG. 2 illustrates a screen shot illustrating an exemplary form view window having a number of fields used to view a single record. As shown, the form view window displays fields for a single record. For example, the header 102 and the body 103, which includes the last name field 106, “Sutton”, first name field 108, “Hal”, events field 110, “30”, and total prize money field 112, “$1,838,740”, are displayed for a single record. Thus, the form view window includes a body defining a plurality of fields, which have a specified order as well as associated attributes. For example, attributes such as degree of rotation, color, and font may be associated with each field. In addition, each of the attributes has a corresponding attribute value. For instance, an attribute such as the font may be specified by values such as Times New Roman. In addition, the order as specified in this example requires that the last name of each of the golf money winners be displayed as the first field of the body 103. Although the form view format provides a simple mechanism to view a single database record, it is often desirable to view and compare multiple database records. For this purpose, the form view report format is inadequate.\nRather than viewing a single record, multiple records are often simultaneously displayed. FIG. 3 is a screen shot illustrating an exemplary list view window used to view multiple records. In addition to the title 102, “Golf Money Winners,” the bodies of multiple records as entered in the main form layout and displayed in the form view format are displayed consecutively. As shown, the body of a first record 302 is shown to include the last name field 304, “Woods”, the first name field 306, “Tiger”, the Events field 308, “20”, and the total prize money field 310, “$1,841,117”. Similarly, the body of a second record 312 is shown to include the last name field 314, “Sutton,” the first name field 316, “Hal,” the events field 318, “30”, and the total prize money field 320, $1,838,740. Although multiple records are displayed, it may be desirable to compare field values of multiple records. For instance, a viewer may wish to compare the total prize money won by multiple golf players. However, each of the records is displayed such that the fields associated with a single one of the records are displayed in multiple rows as well as columns. Moreover, since the records are displayed consecutively in the list view format, the values associated with the same field of multiple records are not displayed adjacent to one another and only a few records at most can be concurrently displayed. As a result, it is difficult to make such a simple comparison between values of the same field. From this list view window, it is difficult to make such a determination. Accordingly, although multiple records may be simultaneously displayed in this manner, the list view format does not facilitate comparisons of fields of multiple records.\nAs described above, it is difficult to compare fields of multiple records using a form view or a list view format. Moreover, it is often desirable to display the same data in a variety of formats. For instance, it may be desirable to vary the format (e.g., appearance) of the fields (e.g., font, column width) as well as the order of the fields. In addition, it may be desirable to alter the manner in which each field is displayed through the association of various attributes such as color, font, or degree of rotation.\nIt is often desirable to display records in a manner suitable for comparing the values of fields of those records. In order to assist a user in creating such a report, database records may be displayed in a format such as a table format (e.g., using formatting information from an existing layout). One method of presenting database records in a table format is disclosed in U.S. Pat. No. 6,613,099, entitled “Process and System for Providing a Table View of a Form Layout for a Database,” listing Christopher Crim as inventor, which is incorporated herein by reference for all purposes.\nReferring now to FIG. 4, a screen shot illustrating an exemplary table view window in which various fields from the form layout may be displayed for one or more records of a database is presented. As shown, the table view window is capable of displaying multiple records 401 such that values associated with the same field of the plurality of records are displayed adjacent to one another. For instance, column header 402 identifies the last name field and the corresponding column includes a column of last names such that the last names associated with the records 401 are displayed adjacent to one another. In this example, the values are displayed in a single column. However, values may similarly be displayed in a single row to facilitate comparison of the values. Values associated with column header 404 identifying the first name field, column header 406 identifying the events field, and column header 408 identifying the total prize money field are similarly displayed for the records 401.\nRegardless of the layout that is used to present database fields, it is often difficult for a user to select from fields in a particular database. For instance, a user may wish to select fields for sorting purposes. Unfortunately, in order to select these fields, the user must typically navigate and interact with the database. For instance, the user must be able to select a database table from numerous database tables in order to select fields from this database table. However, since the database design is set up by a database designer, the user is typically not familiar with the database design. As a result, the user must typically navigate the database to select the desired fields. Unfortunately, since the user is unfamiliar with the database and associated tables, this field selection process is a tedious and time-consuming process.\nIn view of the foregoing, what is needed is a process and system for enabling a user to select fields from a database. Moreover, it would be desirable if the fields could be selected by a user while minimizing the efforts required by the user."}
{"text": "1. Field of the Invention\nAt least one example in accordance with the present invention relates generally to systems and methods for monitoring a load center for current and power usage.\n2. Discussion of Related Art\nA load center or panelboard is a component of an electrical supply system which divides an electrical power feed from a power line into different subsidiary circuit branches. Each subsidiary circuit branch may be connected to a different load. Thus, by dividing the electrical power feed into subsidiary circuit branches, the load center may allow a user to individually control and monitor the current and power usage of each load.\nCurrent Transformers (CT) are commonly used to monitor current in a subsidiary or main branch of a load center. A CT may be used to measure current in a branch by producing a reduced current signal, proportionate to the current in the branch, which may be further manipulated and measured. For example, a CT coupled to a branch of a load center may produce a reduced current AC signal, proportionate to the magnitude of AC current in the branch. The reduced current AC signal may then either be measured directly by measurement instrument or converted to a DC signal and passed to a measurement instrument. Based on the signal received, the measurement instrument may determine the level of current in the subsidiary branch and may assist in providing efficient energy management."}
{"text": "1. Field of the Invention\nThe present invention generally relates to systems for routing telephone calls to appropriate numbers. More particularly, the present invention relates to an Advanced Intelligent Network (AIN) based system and methods for routing telephone calls based on the location of the calling party.\n2. Acronyms\nThe written description provided herein contains acronyms which refer to various communication services and system components. Although known, use of several of these acronyms is not strictly standardized in the art. For purposes of the written description herein, acronyms will be defined as follows:\nAIN--Advanced Intelligent Network PA0 AMA--Automatic Message Accounting PA0 CCIS--Common Channel Interoffice Signaling PA0 CO--Central Office PA0 CPN--Calling Party Number PA0 CPR--Call Processing Record PA0 DN--Dialed Number Trigger PA0 DRS--Data Reporting System PA0 EO--End Office (EO) PA0 ISCP--Integrated Service Control Point PA0 LSP--Local Service Provider PA0 NPA--Number Plan Area, i.e., area code PA0 NXX--Central Office Code PA0 RTN--Routing Telephone Number PA0 SCE--Service Creation Environment PA0 SCP--Service Control Point PA0 SCCP--Signaling Connection Control Part PA0 SMS--Service Management System PA0 SPC--Signaling Point Code PA0 SS7--Signaling System 7 PA0 SSP--Service Switching Point PA0 STP--Signaling Transfer Point PA0 TAT--Terminating Attempt Trigger PA0 TCAP--Transaction Capabilities Applications Protocol\n3. Description of the Related Art\nIn recent years, a number of new telephone service features have been provided by advanced intelligent communications networks such as an Advanced Intelligent Network (AIN). The AIN evolved out of a need to increase the capabilities of the telephone network architecture to meet the growing needs of telephone service customers. The AIN architecture generally comprises two networks, a data messaging network and a trunked communications network. The trunked communications network handles voice and data communications between dispersed network locations, whereas the data messaging network is provided for controlling operations of the trunked communications network.\nAn illustration of the basic components of an AIN architecture is shown in FIG. 1. As shown in FIG. 1, Central Offices (CO) 10-16 are provided for sending and receiving data messages from an Integrated Service Control Point (ISCP) 20 via a Signaling Transfer Point (STP) 30-34. The data messages are communicated to and from the COs 10-16 and the ISCP 20 along a Common Channel Inter-Office Signaling (CCIS) network 22. Each CO 10-16 serves as a network Service Switching Point (SSP) to route telephone calls between a calling station (e.g., station 40) and a called station (e.g., station 48) through the trunked communications network 24-26. For more information regarding AIN, see Berman, Roger K., and Brewster, John H., \"Perspectives on the AIN Architecture,\" IEEE Communications Magazine, February 1992, pp. 27-32, the disclosure of which is expressly incorporated herein by reference in its entirety.\nWhile prior AIN or AIN-type intelligent network applications may have provided various features to subscribers and users, these prior applications do not allow users to dial one telephone number and reach a single point of contact for multiple services provided by a subscriber. Current systems and methods require users to identify one of many possible numbers to call depending on the specific information or service desired from the subscriber. This requires users to know the telephone number of all departments or service groups of the subscriber that they need information from.\nMoreover, none of the current systems and methods allow a user to dial an abbreviated telephone number to access services from a subscriber. Currently, the user must lookup, write down, or memorize a full seven or more digit number for each department or service group that they may need information from.\nTherefore, a system and method is needed that allows users to dial one telephone number and reach a single point of contact for information and services provided by a subscriber, and that provides an abbreviated telephone number that is easy to remember for accessing the single point of contact for services from the subscriber."}
{"text": "Pipes used in offshore oil fields include risers (pipes for pumping up crude oil), umbilicals (integration of pipes for supplying chemicals for crude oil viscosity reduction for the purpose of controlling the pumping, power cables, and others), flow lines (pipes for transporting pumped crude oil which extend on the sea floor), and the like. They have various structures, and known pipes include metal-made pipes and metal/resin hybrid pipes. In order to achieve weight reduction, use of metal-made pipes tends to be reduced and metal/resin hybrid pipes are becoming the mainstream. Since oil drilling sites become much deeper and the temperature of crude oil pumped therefrom rises, resins used for these pipes need to have better mechanical strength and chemical resistance at high temperatures (resistance to high-temperature crude oil, resistance to acidic gas, such as hydrogen sulfide, contained in crude oil at high temperatures, resistance to chemicals such as methanol, CO2, and hydrogen chloride injected so as to reduce the crude oil viscosity at high temperatures), and lower permeability at high temperatures. Thus, there is a demand for materials which can take the place of polyamide (the operating temperature range is up to 90° C.) and polyvinylidene fluoride (the operating temperature range is up to 130° C.) which have been used for the pipes.\nPatent Literature 1 discloses as a material suitable for flexible pipes a fluororesin which is a copolymer containing copolymerized units of tetrafluoroethylene, vinylidene fluoride, and an ethylenic unsaturated monomer excluding tetrafluoroethylene and vinylidene fluoride, and has a specific storage elastic modulus."}
{"text": "1. Technical Field\nThe present invention is directed generally to wireless communication systems and, more particularly, to a system and method for parameter selection to avoid interference in a wireless communication system.\n2. Description of the Related Art\nEarly wireless communication devices, commonly known as cell phones, provided wireless voice services to the user. These early phones have been replaced with wireless communication devices capable of delivering voice, data, and multi-media information. In addition, wireless devices often include location determination using the Global Positioning System (GPS). The delivery of these additional services requires additional bandwidth. In some cases, bandwidth previously allocated for one purpose has been reassigned for the delivery of wireless communication services. For example, the spectrum originally allocated to Ultra-High Frequency (UHF) television has been partially reallocated for wireless communication services.\nDevices are being designed with multiple services that depend on multiple radio systems being operated at the same time. For example, devices are being designed that can connect to the cellular network using several different radio protocols and frequency bands. In addition, these devices may have other applications, such as broadcast television or Bluetooth, which use independent radio systems.\nThese independent radio systems may interfere with, or be interfered by, the radio system used for cellular operation. One can appreciate that the operation of multiple transceivers within a single device may decrease the operational capability of the device. Therefore, it can be appreciated that there is a significant need to reduce interference among the multiple transceiver systems. The present invention provides this, and other advantages, as will be apparent from the following detailed description and accompanying figures."}
{"text": "FIELD OF THE INVENTION\nThis invention pertains to increasing the number of building blocks that can be incorporated independently into oligonucleotides via DNA and RNA polymerases, and nucleoside analogs capable of forming non-standard Watson-Crick base pairs joined by patterns of hydrogen bonding different from those found in the adenine-thymine and cytosine-guanine base pairs."}
{"text": "1. Field of the Invention\nThe present invention relates to the art of adhering a thin film to a base plate, and particularly relates to the art in which a dry film consisting of a light-transmissible resin film and a photosensitive resin layer provided on one side of the resin film is stuck under pressure to a base plate such as a printed circuit board and an electronic circuit substrate of silicon, gallium arsenide or the like without leaving an air bubble between the dry film and the base plate and without wrinkling the dry film.\n2. Description of the Prior Art\nIn one conventional apparatus for adhering a dry film to a base plate under pressure, as disclosed in Japanese Patent Publication No. 13341/80, the base plate and a photosensitive resin layer bonded to one side of a light-transmissible resin film are brought into contact with each other while being kept in a vacuum chamber of reduced pressure. Sufficient pressure is thereafter applied to the other side of the resin film to push the film and the layer onto the base plate. Thereafter, heat is applied. In another conventional apparatus the dry film is pressed onto the base plate in the air at atmospheric pressure by a rotating heat and pressure sticking roller. In yet another conventional apparatus as disclosed in the Japanese Patent Application (OPI) No. 121696/83, the end of the base plate and that of the dry film are stuck to each other in the air at atmospheric pressure and the film and the plate are thereafter completely stuck to each other in a vacuum. The dry film consisting of a light-transmissible resin film and a photosensitive resin layer bonded to one side of the film is thus stuck to the base plate, under prescribed pressure and at a prescribed temperature in each of the conventional apparatuses. The photosensitive resin layer is thereafter exposed to light through a pattern mask overlaid on the light-transmissible resin film. The resin film is then removed. The resin layer is then etched so that a desired pattern is made on the base plate.\nIn the above-mentioned conventional arts, however, an air bubble is likely to be left between the dry film and the base plate at the time the film and the plate are stuck to each other. The air bubble causes defects in the step of exposing the film and layer to light or the step of etching the resin layer, thereby causing problems for the process. If the dry film is entirely or partly stuck to the base plate in air at atmospheric pressure, not only will air bubbles likely be left between the dry film and the base plate but also the film will likely be wrinkled by the pressure of the rotating pressure sticking roller. This is another problem."}
{"text": "This invention relates generally to the cancellation of an echo signal in a voice communication system.\nIn worldwide telecommunications systems, echoes arise in various situations and impair communication quality. Echoes occur when a delayed or distorted version of an original audio signal is reflected back to the source. In telecommunications networks, an impedance mismatch is one factor that contributes to the refection of an audio signal back to its source. The reflected audio signal is a delayed or distorted version of the original signal, which causes echoes in speech communication systems.\nA hybrid transformer is typically used to connect a two-wire local telephone exchange to a four-wire long distance or mobile telephone network. The imperfect impedance match exhibited by the hybrid transformer generates the echoed signal. In the past two decades several methods have been used to alleviate the echo problem and improve communication quality. These prior art methods are collectively referred to as echo cancellation.\nThe long distance or mobile telephone from which a voice signal originates is commonly referred to as the xe2x80x9cfar-endxe2x80x9d. The voice signal from the far-end is called the inbound signal and travels through a path called the receive-path. The inbound signal passes through a hybrid transformer located at a local telephone exchange. The hybrid transformer is typically made integral to a device called a Central Office Line Interface Unit. Most of the inbound signal is transferred through the hybrid transformer to the party subscribing to the local telephone exchange that is receiving the phone call. The subscriber using the local telephone exchange is referred to as the xe2x80x9cnear-endxe2x80x9d. The hybrid transformer propagates a signal originating at the near-end, commonly called the xe2x80x9cnear-end signalxe2x80x9d, to the far-end using a second signal path called the xe2x80x9csend-pathxe2x80x9d. An unwanted version of the inbound signal is also coupled into the send path resulting in an echo. This unwanted version of the inbound signal is the echo that needs to be eliminated. The composite of the near-end signal and the reflected inbound signal is referred to as the xe2x80x9coutbound signalxe2x80x9d.\nThe echo-path-model is a transfer function that describes the amount of the inbound signal that is reflected back into the outbound signal. In order to determine the echo-path-model, echo cancellation systems monitor the inbound signal and compare that inbound signal to the amount of echo signal observed in the send-path. This process can only be accomplished when the send-path is devoid of any other signals. When the near-end is generating a signal, the presence of that near-end signal in the send-path will preclude any estimation of the echo-path-model.\nOnce the echo-path-model has been derived, an estimate of the echo signal can be calculated. The estimated echo is subtracted from the send-path leaving only the desired near-end signal. Because the resulting transfer function for the echo-path-model is only an estimate of the actual echo transfer function, some of the echo signal will remain in the send-path. This component is called the residual echo.\nEcho cancellers usually use some form of filter to implement the echo-path-model. By subjecting the inbound signal to the filter, an estimate of the echo can be derived. The filter itself is normally an adaptive filter that can be based on one of many different adaptation algorithms. One such algorithm is the Least Mean Squares (LMS) algorithm. To support an LMS based implementation of an echo canceller, a coefficient generator is used to sample both the inbound signal and the outbound signal. From these two signals, a set of filter coefficients are determined and fed to the LMS filter. Again, it is important to note that the coefficient generator cannot perform its function if there is a near-end signal present in the send-path.\nAs the echo canceller continues to operate, the residual echo is used to adjust the coefficients of the LMS filter that models the echo-path. This process is called adaptation. As the adaptation process continues, the coefficients of the filter assume values that more accurately represent the actual echo-path-model. When the coefficients of the filter no longer change, the echo canceller is said to have converged and a near-perfect echo estimate can be derived.\nBecause the outbound signal is a composite of the reflected component of the inbound signal and the near-end signal, it is impossible to measure the magnitude of the reflected echo signal in the presence of the near-end signal. To overcome this, echo cancellation systems normally comprise a double-talk detector that senses when the near-end signal is active. The double-talk detector sends a signal to the coefficient generator that causes the coefficient generator to suspend the adaptation process.\nThe actual echo-path in any given system constantly changes as a result of varying physical phenomenon experienced by the system components themselves. Because of these variations, the adaptation process will seldom converge in a perfect echo-path-model.\nOne way to improve the accuracy of the echo-path-model is to ensure that the adaptation process is performed as quickly so that any temporal variations in the signal line can be reflected in the resulting filter coefficients. By achieving faster convergence, echo-cancellation systems could reduce the amount of residual echo remaining in the send-path. This would contribute to better voice quality in the communications system.\nIn one illustrative embodiment of the present invention, an estimate of the spectral distribution of an inbound signal is used as a basis for filter coefficients for a filter disposed prior to a coefficient generator and an echo-estimation filter. This first filter flattens, or whitens the spectrum of the inbound signal used to generate coefficients and is likewise subjected to the echo-estimation filter. The echo-estimation filter actually implements the transfer function for an echo-path-model that describes the system.\nIn this same illustrative embodiment, a second filter is placed in send-path prior to a subtractor that is used to subtract an estimated echo from the send-path. This second filter uses the same coefficients used by the first filter. The second filter flattens the spectrum of the outbound signal. Hence, the adaptation filter operates on spectrally equalized versions of the inbound and outbound signals. Once the estimated echo is subtracted from the send-path, the outbound signal is fed through a reconstruction filter in order to introduce the original spectral components of the inbound signal into the equalized outbound signal. By flattening the inbound and the outbound signal, the adaptation filter will converge to a solution of an echo-path-model in less time compared to conventional echo cancellation systems. This contributes to better overall echo cancellation quality.\nThere are, of course, several brute force mechanisms for achieving faster convergence in an echo cancellation system. These brute force methods rely principally on the use of faster processors in the implementation of the adaptive filters. The present invention exploits the fact that certain adaptive filters converge more rapidly when the input signal presented to the filter has been equalized.\nThe present invention comprises in the first instance a method for canceling echoes in communications systems. When an inbound signal is received, the method provides that the frequency spectrum of the inbound signal should be determined. Determining the spectrum of the inbound signal can be accomplished several ways. In one example embodiment, the inbound signal is actually measured and the spectrum is determined from the measurement. In an alternative embodiment, the general characteristics of the communications system are monitored over some period of time. Based on the historical observations of the communications systems channel, an exemplary spectrum is determined and subsequently used in the process. Once the frequency spectrum of the inbound signal is determined, the inbound signal is itself equalized.\nIn this example embodiment of the method for canceling echoes, an outbound communications signal is also equalized based on the frequency spectrum of the inbound signal. The outbound communications signal typically comprises at least two components. These are the actual near-and signal that must be propagated to a far-end and some derivative of the inbound communications signal; i.e. the echo component. Using the flattened inbound communications signal and the flattened outbound communications signal, filter coefficients are generated for an adaptive filter. The adaptive filter is a convenient means for implementing the echo-path-model that the communications system exhibits. As such, subjecting the inbound communications signal to the adaptive filter results in an estimate of the echo component found in the outbound communications signal.\nOnce an estimate of the echo component is generated by the adaptive filter, the method provides for subtracting the estimated echo component from the flattened outbound communications signal. In theory, the outbound communications signal should be devoid of any echo component at this stage. Prior to directing the echo-canceled outbound signal to the far-end, the signal must be reconstructed so as to include the original spectral envelope representative of the original inbound signal. This reconstruction is accomplished based on the spectral distribution for the original inbound signal.\nIn one example embodiment, the method of the present invention provides for monitoring the frequency spectrum of the inbound signal by buffering the inbound communications signal and then calculating correlation coefficients for the buffered signal. The correlation coefficients are used to create a set of normal equations that can then in turn be solved to determine the frequency spectrum of the inbound signal.\nThe invention further comprises a system for canceling echoes. In one example embodiment, an echo cancellation system comprises a receiver capable of excepting external signal and then propagating that signal to other components in the system. The echo cancellation system further comprises a spectrum estimator that is able to create filter coefficients. These filter coefficients are used to configure a first whitening filter that accepts the inbound signal from the receiver and creates an equalized rendition of the inbound signal.\nThe receiver also propagates the inbound signal to a hybrid transformer. The hybrid transformer, having received the non-equalized inbound signal, directs the non-equalized inbound signal to a near-end subscriber. Unfortunately, not all of the inbound signal is propagated to the subscriber. Some of the inbound signal is coupled together with a near-end signal that originates with the subscriber. The inbound signal that is coupled together with the near-end signal is the echo signal that needs to be cancelled. The combination of the near-end signal and the echo signal is called the outbound signal. The invention further comprises a second whitening filter that receives the outbound signal and equalizes the outbound signal based on the filter coefficients created by the spectrum estimator.\nThe flattened inbound signal and the flattened outbound signal are initially used to create filter coefficients for an adaptive filter that further comprises the invention. The adaptive filter, as configured by these filter coefficients, implements the echo-path-model transfer function that can be used to estimate the amount of echo that should be found in the outbound signal. The invention further comprises a subcontractor that subtracts the predicted echo from the outbound signal.\nBecause the outbound signal has been equalized, it must be reconstructed in order to reflect the spectral envelope of the original inbound signal. This is accomplished by a reconstruction filter that further comprises the invention and whose filter coefficients may be based on the inverse of the inbound signal\"\"s spectral density.\nIn one example variant of the present invention, the spectrum estimator comprises a spectrum analyzer that actually measures the inbound signal and determines the level of energy at various frequencies. Such a spectrum analyzer may comprise a buffer that captures the inbound communications signal and a correlation calculator that creates correlation coefficients based on the inbound signal stored in the buffer. In this example embodiment, the equation generator creates a set of normal equations based on the coefficients created by the correlation calculator. A matrix analyzer solves this set of normal equations; the resulting matrix defines the spectral distribution of the inbound signal.\nIn yet a second illustrative variation of the present invention, the spectrum estimator comprises a memory for storing a set of anticipated spectral distributions. In this example embodiment, the invention further comprises a selection unit that selects the estimated spectrum from the memory. The anticipated spectral distributions may be created off-line and can be based on a priori knowledge of the condition of the inbound communications line or can be based on an extrapolation of historical observations of line condition.\nIn all of these example embodiments, any of the whitening filters can be implemented as filter processors that accepts coefficients from the spectrum estimator. The reconstruction filter may also comprise a filter processor that accepts filter coefficients from the spectrum estimator.\nThe present invention further comprises a central office line interface that can be used in a telephone-switching center or like application. In one illustrative embodiment, the central office line interface unit according to the present invention comprises a hybrid transformer, the first whitening filter, a second whitening filter, a coefficient generator, an adaptive filter, a subtractor, and a reconstruction filter. The central office line interface unit may further comprise a double-talk detector.\nIn this example embodiment, the central office line interface unit receives the inbound signal from a remote exchange using a four-wire interface provided by the hybrid transformer. The hybrid transformer directs the inbound signal to a two-wire interface that is used to service a local subscriber. The four-wire interface provided by the hybrid transformer itself comprises a two-wire send-path and a two-wire receive-path. In operation, a near-end signal is received from the local subscriber and is directed into the two-wire send-path by the hybrid transformer.\nThe inbound signal arrives at the central office line interface unit by way of the two-wire receive-path provided by the hybrid transformer. The first whitening filter equalizes the inbound signal and creates a flattened inbound signal. The second whitening filter concurrently flattens the outbound signal emanating from the hybrid transformer on the two-wire send-path. It should be noted that this outbound signal comprises a near-end signal originating with the local subscriber and an echo signal that is coupled into the send-path by the hybrid transformer.\nThe coefficient generator receives the flattened inbound signal and the flattened outbound signal and creates filter coefficients for the adaptive filter. The adaptive filter, as configured by these filter coefficients, implements the echo-path-model transfer function that is used to predict the nature and quality of the echo signal coupled into the send-path by the hybrid transformer. In operation, the adaptive filter will create the estimated echo signal. The subcontractor receives the flattened outbound signal from the second whitening filter and subtracts the estimated echo signal therefrom.\nIn another example embodiment of the present invention, the coefficient generator may continually refine the coefficients that define the echo-path-model in a process called adaptation. This is done by receiving a residual echo signal from the subcontractor. Such adaptation, and for that matter the creation of the original echo-path-model coefficients can only be accomplished in the absence of any near-end signal. This is due to the fact that the near-end signal obscures the echo signal necessary to determine the echo-path-model. The central office line interface unit of the present invention may in this instance further comprise a double-talk detector. The double-talk detector monitors the state of the near-end signal source. When the near-end signal is active, the double-talk detector issues a signal that prevents the coefficient generator from updating its coefficients.\nIn yet another example embodiment of the present invention, the first and second whitening filters and the reconstruction filter are configured based on the estimate of the spectral distribution of the inbound signal. The estimate of the spectral distribution of the inbound signal may be obtained either through measurement of the inbound signal and a determination of the spectral distribution thereof or by simply anticipating what the spectral distribution of the inbound signal would be under some given circumstance."}
{"text": "This application relates to optical resonators, and more specifically, to optical whispering-gallery-mode resonators.\nA dielectric sphere may be used to form an optical whispering-gallery-mode resonator which supports a special set of resonator modes known as xe2x80x9cwhispering gallery modesxe2x80x9d. These modes represent optical fields confined in an interior region close to the surface of the sphere around its equator due to the total internal reflection at the sphere boundary. Microspheres with diameters on the order of 10xcx9c102 microns have been used to form compact optical resonators. Such resonators have a resonator dimension much larger than the wavelength of light so that the optical loss due to the finite curvature of the resonators can be small. The primary sources for optical loss include optical absorption in the dielectric material and optical scattering due to the inhomogeneity of the sphere (e.g., irregularities on the sphere surface). As a result, a high quality factor, Q, may be achieved in such resonators. Some microspheres with sub-millimeter dimensions have been demonstrated to exhibit very high quality factors for light waves, exceeding 109 for quartz microspheres. Hence, optical energy, once coupled into a whispering gallery mode, can circulate at or near the sphere equator with a long photon life time.\nWhispering-gallery-mode resonators may use resonator geometries based on spheres. Since the whispering gallery modes essentially exist near the equator of a sphere, a resonator may not be necessarily a whole sphere but a portion of the sphere near the equator that is sufficiently large to support the whispering gallery modes. Hence, rings, disks and other geometries formed from a proper section of a sphere may be used. Such resonators are still spherical and their whispering gallery modes are essentially identical to such modes of the respective whole spheres.\nThis application includes non-spherical whispering-gallery-mode resonators and their applications. In one embodiment, a spheroidal cavity is used to generate a true finesse on the order of 104, a free spectral range of the order of few nanometers, and a quality-factor Qxcx9c1xc3x97107. Devices based on such a spheroidal cavity are also disclosed."}
{"text": "1. Field of the Invention\nThe present invention relates to a closure for body fluid collection, transport or storage containers and, more particularly, relates to a ball and socket closure to be used to resealably close a container being used in a laboratory or other clinical environment.\n2. Background Description\nAfter a doctor, phiebotomist or nurse has used an evacuated blood collection tube or other primary tube to draw a primary sample of body fluid from a patient in a hospital or doctor\"\"s office, the primary sample will typically be xe2x80x9cpoured offxe2x80x9d or pipetted into a secondary tube so that the sample can be simultaneously tested in two or more different areas of a clinical chemistry laboratory. For example, the sample may undergo routine chemistry, hormone, immunoassay, or special chemistry testing. In addition, the sample is sometimes xe2x80x9cpoured off or pipettedxe2x80x9d into a secondary tube for overnight storage, to transport the sample from one laboratory to another, or to remove the plasma or serum sample from a separator gel or red blood cells used in the primary tube. When the secondary tube is not being used or is being transported, it is very important to close the open end of the secondary tube with a closure to prevent contamination, evaporation or loss of the sample.\nCurrent closures for secondary tubes include plastic caps that snap over or into the secondary tube or cork or rubber stoppers, wherein the stopper is solid and includes a plug portion that fits in the open end of the tube and an enlarged head portion used to remove the closure from the tube using a two-handed method. Such closures provide means for sealing the open end of the tube, but are difficult to remove with two hands and impossible to remove using only one hand. This presents a problem, since the closure must be removed from the tube and discarded prior to placing the tube in a chemical analyzer due to the inability of most sample probes to penetrate any solid closure material. In view of the above, it is desirable to have a closure that can be easily removed from the tube or a closure that can remain on the tube and be easily opened and closed many times for manual sample access and/or during dire sampling by a chemical analyzer.\nThe present invention overcomes the problems identified in the background material by providing a closure for primary or secondary fluid collection, transport or storage containers or tubes for body fluids that can easily be opened and closed multiple times.\nA preferred embodiment of a closure according to the present invention includes a bail and socket closure to be used to resealably close a specimen container or s tube used in a laboratory or other clinical environment. In one embodiment, the ball and socket closure is snap-fitted into a tube. The ball has a tab extending therefrom that is pushed by a user approximately 90 degrees to rotate the ball within the socket to a position wherein a passageway through the ball aligns with the opening of the tube and provides access through the closure to the inside of the tube. When the tab is pushed 90 degrees in the opposite direction the ball rotates to close the passageway and seal the open end of the tube for storage to avoid evaporation and for possible access or retest at a later date.\nAn object of the ball and socket closure of the present invention is to provide dirt access to the tube such that a transfer pipette or an analyzer sample probe can access the fluid contents of the tube without the probe contacting the inner surface of the tube or the closure itself. This structure prevents contact or contamination of the probe while maintaining a one handed closure operation. The tab on the ball provides for an easy opening and closing operation with one hand during use which is also a major ergonomic and workflow improvement over existing closures and tubes.\nAnother object of the present invention is to provide a closure having an outer diameter that is no larger than the outer diameter of a current primary specimen collection container with closure (i.e., the VACUTAINER(copyright) SST(copyright) Brand Tube sold by Becton Dickinson and Company) so that the entire closure and tube assembly can be loaded into conventional analyzer racks, carousels or holders without removing the closure from the tube. Since the closure does not need to be removed from the tube, risk of loss or accidental contamination is minimized.\nIn addition, the ability to use only one closure through multiple samplings rather than replacement of stoppers multiple times reduces cost for the user.\nIn addition, the closure of the present invention is dimensioned to develop a liquid seal that prevents any liquid from leaking out of the tube through or past the ball and socket closure when it is in the closed position.\nThese and other aspects, features and advantages of the present invention will become apparent from the following detailed description taken in conjunction with the accompanying drawings."}
{"text": "In a code division multiple access mobile communication system, since many channels use the same frequency, the reception power of a signal on a given channel becomes interference wave power that interferes with other channels. On an upstream channel through which a mobile station transmits a signal and a base station receives it, when the ratio of signal power to interference wave power becomes high and excessive reception quality is set, the interference wave power increases. As a consequence, the channel capacity decreases. In order to prevent this, the transmission power on the mobile station must be strictly controlled. Transmission power control on an upstream channel is performed in the following manner. A base station measures reception quality based on a signal-to-interference ratio or the like, and compares it with a control target value. If the reception quality is higher than the control target value, the base station transmits a control instruction to decrease the transmission power to the mobile station. If the reception quality is lower than the control target value, the base station transmits a control instruction to increase the transmission power to the mobile station. The mobile station then increases/decreases the transmission power in accordance with the control instruction. This transmission power control method is disclosed in detail in U.S. Pat. No. 5,056,109 (Gilhousen et al., “Method and apparatus for controlling transmission power in a CDMA cellular mobile telephone system”).\nOn a downstream channel as well, a high channel capacity is realized by performing transmission power control to set reception quality based on a signal-to-interference ratio or the like to a predetermined control target value. In transmission power control on a downstream channel, a mobile station measures the reception quality of the downstream channel, and compares it with a control target value. If the reception quality is higher than the control target value, the mobile station transmits a control instruction to decrease the transmission power to the base station. If the reception quality is lower than the control target value, the mobile station transmits a control instruction to increase the transmission power to the base station. The base station then increases/decreases the transmission power in accordance with the control instruction.\nIn an actual propagation environment for a mobile communication system, since mobile stations differ in their multipath effects and moving speeds, if a constant control target value is set for the above transmission power control on upstream and downstream channels, channel quality based on a bit error rate, frame error rate, or the like cannot be kept constant. If a large control target value is uniformly set to satisfy predetermined channel quality in mobile stations in any conditions, the set control target value is unnecessarily large for many mobile stations. As a consequence, the transmission power becomes unnecessarily high accordingly to increase interference wave power that interferes with other channels, resulting in a reduction in channel capacity. Optimal control target values that can minimize an increase in interference wave power and obtain a predetermined channel quality differ from mobile station to mobile station.\nAs a method of controlling a control target value to an optimal value, a technique called an outer loop is available, which changes a control target value in accordance with channel quality. When a frame error rate is to be used as channel quality, an error detection code is set in each frame, and a frame in which an error is detected by this code is determined as an error. A control target value is changed to set the frame error rate to a predetermined channel quality target value.\nA specific method for this operation is described in Higuchi, Ando, Okawa, Sawabashi, and Adachi, “Experimental Study on Adaptive Transmission Power Control Using Outer Loop in W-CDMA” (Technical Report of IEICE, RCS98-18 (1998-04), pp. 51–57). As described in this reference, if the frame error rate in a predetermined time is higher than a target frame error rate, the control target value is increased by a predetermined amount, whereas if the frame error rate in the predetermined time is lower than the target frame error rate, the control target value is decreased by the same predetermined amount. The frame error rate in a predetermined time is obtained by counting the number of frames determined as errors for each predetermined number of frames and dividing the number of frames determined as errors by a predetermined number.\nAssume that, in this method, a low frame error rate is set as a channel quality target value, and a control target value is changed by using the frame error rate in a predetermined time which is obtained from a small number of frames. In this case, even with the same channel quality, since the frame error rate in the predetermined time varies, a control target value is often set apart from an optimal value. Assume that in order to prevent this, the control target value is changed by obtaining a frame error rate within the predetermined time from many frames. In this case, when the optimal control target value changes, it takes much time to change the control target value in accordance with this change. Assume that the moving speed abruptly changes, the optimal control target value increases, and many frame errors are caused. Even in this case, the control target value cannot be increased until a predetermined number of frames are received, and the frame error rate in the predetermined time is calculated. As a consequence, the state where many frame errors are caused will continue. In contrast to this, if the optimal control target value decreases, the control target value cannot be decreased until a predetermined number of frames are received. As a consequence, the state where the transmission power is unnecessarily high will continue, resulting in a reduction in channel capacity.\nAnother problem in this conventional method is that the frame error rate greatly varies over time. In a constant propagation environment, even if the frame error rate is fixed to an ideal control target value as a channel quality target value, the frame error rate within a short time unit varies depending on the number of frame errors caused in the time. According to the conventional method, however, even if a propagation environment such as a moving speed is constant, since a control target value is repeatedly increased and decreased, the frame error rate within a short time unit greatly varies as compared with a case where a control target value is fixed to an ideal control target value. This is because, when the control target value is set to be larger than an optimal value, the frame error rate becomes lower than the channel quality target value, whereas when the control target value is set to lower than the optimal value, the frame error rate becomes higher than the channel quality target value.\nWhen voice and image are to be transmitted in real time, even the frame error rate in a short time unit must be suppressed to be equal to or lower than a predetermined value because the quality of the voice and image deteriorates if frame errors concentrate. Therefore, a large channel quality target value must be set to make the frame error rate fall within a predetermined value while the frame error rate increases. As the control target value is increased, the transmission power increases, resulting in a reduction in channel capacity."}
{"text": "The present invention relates to an apparatus for reproducing video data and audio data, and more particularly to an apparatus for reproducing data and the apparatus for recording data which are arranged to reproduce and/or record data from/onto a recording medium based on copy management information.\nA DVD-ROM has about seven times as large a volume as a CD-ROM. The DVD-ROM may contain program codes for PCs as well as motion picture software created by compressing video data and audio data. The DVD also includes a DVD-RAM, a DVD-R, and a DVD-RW as its recording medium. These recording medium may record a large amount of data, which brings about a possibility of illegally copying software such as a motion picture. This sort of illegal copy has to be prevented. Hence, a technique of preventing illegal copy becomes important. This technique is described in the magazine “Nikkei Electronics” issued on Aug. 18, 1997, pages 110 to 119, published by Nikkei BP publishing.\nFor example, the motion picture recorded on the DVD-Video disk is typically coded according to the CSS (content scrambling system). Hence, the copied data cannot be reproduced unless it is descrambled."}
{"text": "Battery is the electric energy source of terminal devices, and if there is no battery, the application performance of the terminal device will be greatly reduced. A high-capacity high-performance battery can not only provide long duration capability for a terminal but can also protect the terminal, so that the terminal device can operate for a long time with high efficiency.\nWith the popularization of modern mobile phones, especially smart phones, the functions thereof are richer and richer, which have become necessary for user communication, office work and entertainment. Most people think highly of battery duration capability, 38% interviewees state that as long as the mobile phone can continue operating normally (such as dialing and answering, and receiving and sending short messages), they are satisfied and do not want to sacrifice the battery duration time for any additional functions. The data obtained from the after-sales market also demonstrates that customers do not satisfy with the service life of the battery being shorter and shorter and the charging efficiency being lower and lower to the greatest extent.\nThe reason that the service life of the battery becomes shorter and shorter may be: frequent and repeated high-potential charging and over-charging, battery over-discharging due to not be charged for a long time, etc. For common users, from the moment a terminal device is used, the performance of the battery will decrease gradually over time, and therefore, to retain the performance of the battery to the greatest extent means a lot."}
{"text": "The present invention pertains to an environmentally controlled building. More particularly, the present invention pertains to a building providing a healthful environment beneficial for people living, working, or spending leisure time in the building, while avoiding temperature extremes, polluted air, sudden variations in barometric pressure, and other conditions which are detrimental to the health and well being of the inhabitants. The building encloses a covered atrium or courtyard which can have an air supported continuous membrane fabric roof or a conventional roof.\nMany people desire or require healthful, controlled environments in which to live, work, and engage in leisure activities. This is particularly true of older persons and of people having health problems such as respiratory problems, arthritis, or rheumatism. It is a common practice to control the temperature and humidity of the air within a building. More than simple temperature and humidity control are desirable, however. Thus, it is also desirable to remove pollen and other sources of pollution from the air. Additionally, sudden changes in ambient atmospheric pressure can have an adverse effect on people, particularly people bothered by arthritis or rheumatism, and so atmospheric pressure changes should be controlled."}
{"text": "The invention relates to a method and apparatus for distribution and mixing of high concentration or consistency (5% or higher) cellulose pulp with a treatment fluid, such as chlorine or chlorine dioxide.\nThe object of the invention is to make such distribution and mixing as effective as possible, so that the treatment fluid(s) is distributed as evenly as possible in the pulp suspension when introduced thereto, so that mixing of the pulp suspension and treatment fluid is effected, so that even a relatively small quantity of a treatment fluid is distributed evenly in and around all particles or fibers of the pulp suspension.\nThe effectiveness of such distribution and mixing depends on many factors, such as the pulp concentration in relation to the quantity of liquid or gas which is to be added, the solubility of the added liquid or gas in the suspension liquid, and to the reaction speed of the added treatment fluid with the particles of the pulp suspension. Generally, it can be said that the higher the concentration of solids or fibers in the pulp suspension, the more difficult it is to mix in treatment fluids so that they are evenly distributed in the suspension. Generally, it can also be said that the faster the added fluids react with the pulp, the more important it is that the fluids are distributed and mixed in as quickly and as evenly as possible. Since chlorine reacts quickly with pulp, and since it is desirable to treat high solids concentration pulp during bleaching, it is especially important to quickly mix chlorine with pulp. Since chlorine has an especially quick initial reaction with pulp and since it is undesirable to dilute the pulp with additional quantity of liquid, chlorine is most often added as gas dispersed in a relatively small quantity of liquid which, however, in turn means that problems can easily arise in the distribution and mixing of such a relatively small quantity. An object with the invention is therefore to solve this problem and also to solve the problems which arise when the pulp suspension has relatively high consistency of fibers, preferably above 5%, e.g., about 8-20% or about 10%.\nIn the pulp industry bleaching of pulp with chlorine liquid has hitherto preferably been done at 3-4% concentration mainly due to mechanical difficulties with mixing in and distribution [gas phase chlorination may be done with a pulp concentration in the range of 20-50%]. Since in other treatment stages of industrial bleach plants the pulp concentration normally is kept around 10%, it is desirable also to be able to effect the treatment of pulp with chlorine at this same concentration so that one can use uniform equipment in the bleach plant. This has special importance for the washing apparatus which is used between the treatment stages. Since the treatment with chlorine most often takes place in the beginning of the bleach plant and the pulp therefore must be thickened to about 10% concentration before the pulp goes on to the next treatment state, simplification and bulk reduction of equipment can be obtained if this first chlorine treatment also can take place at about the same high concentration.\nThe present invention allows chlorine treatment with proper mixing of high concentration pulp, as will become clear from an inspection of the detailed description of the invention and the appended claims."}
{"text": "Metal cans have been known for having a structural top flange wall secured to the upper edge of the can tubular body and including a depending inner skirt defining the opening for accessing the inside of the can. The skirt also serves as a seat, to which is pressure fittable the peripheral wall of a sealing lid. The lid is manually removable and also recloseable during the period the can is used.\nIn these prior art cans, the sealing and the axial retention of the lid in the mounted position are achieved by pressure seating a peripheral wall of the lid against the inner surface of the depending skirt of the can flange wall. This type of construction has some deficiencies resulting from the small amount of axial locking force of the lid to the can and also from the fact that the free lower edge of the skirt is in contact with the product in the can.\nIn order to eliminate the aforementioned deficiencies, there has been developed a can and lid as described in the copending patent application P19600454 currently to U.S. Pat. No. 5,899,352 of the same applicant, according to which the axial locking of the lid is achieved by providing, in one of the parts defined by the can skirt wall forming the opening and the lid peripheral wall of an annular rib that is fittable in a corresponding circumferential recess on the lid peripheral wall or can skirt and having a section similar to that of the annular rib. In this construction, the rib and the recess have substantially coincident profiles of a semicircular shape, defined so that said parts fit each other with a substantially complete contact between the confronting surfaces.\nWhile this arrangement provides a substantial axial force for locking of the lid, avoiding its undue opening due to shocks, internal pressure increase, etc., and allowing an adequate degree of sealing of the can contents to be achieved, keeping the canned product out of contact with air, this construction has the disadvantage of requiring great precision in the formation of the annular rib and circumferential recess.\nDue to dimensional imperfections that can exist in these cans, the fitting between the rib and the corresponding circumferential recess sometimes can have a radial gap, thereby reducing the contact between the respective confronting surfaces to only one point of tangency along a line that develops around the circumference of the rib and recess. This contact does not guarantee an adequate sealing for the can, allowing the canned product to deteriorate.\nAnother deficiency of the known constructions for the can and lid refers to the achievement of automatic closing of cans at the filling units. Cans in which the conventionally secured peripheral edge of the inner ring of the can is on a plane which is at the same level or slightly above the plane of the upper edge of the lid seated at the central opening in some cases have an inadequate closing of the lid. There will be insufficient introduction and pressure of the lid against the can when the latter is moved under the closure roll or piston of the filling machine.\nStill another deficiency of the known constructions refers to the accumulation of the product which is spilled over the structural ring of the can during the progressive removal of the can contents, making subsequent closings of the lid more difficult and consequently making possible the contact of air with the product inside the can. This allows, for example, the oxidation of the product and also, in the case of products having volatile elements in their composition, such as paints, the evaporation of said elements, causing the hardening of the remaining product in the can."}
{"text": "It is known to utilize direct electron lithography and photolithography during respectively different steps of a process for fabricating a microminiature integrated device. In such a process, an electron beam exposure system is advantageously employed to define some of the more critical features of the device. The other features are defined photolithographically.\nFor the electron lithographic step(s) of such a hybrid process, highly sensitive electron resists are available. By utilizing these resists, it is economical in some cases to expose even large areas of a resist-coated wafer with an electron beam system. But, in practice, such resists are typically characterized by (1) relatively poor resolution of developed patterns in thick films, (2) relatively poor tolerance to many dry etching processes of practical importance and (3) the disadvantage that the substitution of electron resists for photoresists in a photolithographic fabrication sequence requires modification of a number of the standard photolithographic processing steps other than the exposure step itself. For these reasons in particular, proposals to utilize an electron beam system to complement a photolithographic device fabrication process have not heretofore usually appeared attractive.\nMoreover, in such a hybrid fabrication process, it appeared not to be feasible to expose a relatively insensitive photoresist (rather than a sensitive electron resist) with a high-speed electron beam system of the raster scanning type."}
{"text": "U.S. Pat. No. 6,656,986 teaches various polyethylene, peroxide-crosslinkable compositions useful in the manufacture of power cable insulation. Some of these compositions have achieved commercial success in the medium voltage power cable market, and an interest exists in extending these commercially successful compositions into the high and extra high voltage power cable markets.\nThe manufacture of power cable insulation is a multistep process that can be separated into two broad parts, i.e., first making a composition from which the cable insulation is made, and second, extruding the composition over single or stranded conductor as an insulation.\nIn one embodiment of the first part of the process, i.e., the part in which the composition is made, a base polymer, e.g., polyethylene, is mixed with one or more additives and then formed into pellets which are soaked with peroxide and subsequently stored and/or shipped to a fabricator who performs the second part of the process, i.e., converting the pellets to a wire or cable coating. To avoid acid catalyzed decomposition of peroxide during storage and shipping, U.S. Pat. No. 6,656,986 teaches inclusion of oligomeric and/or high molecular weight hindered amine stabilizers (HAS).\nIn the making of the pellets, care is taken not to introduce or create impurities that can adversely affect the utility of the composition once formed into a wire or cable sheath. However, some impurities are inevitably introduced into the composition either as, for example, contaminants associated with feed materials to the process, or are made during the process as, for example, gels that result from degradation of the base polymer. Efforts are made, of course, to minimize and remove these impurities before the composition is extruded as a power cable sheath. Some of the impurities are in the form of fine, e.g., less than 100 microns (μm), particulates and are susceptible to removal from the composition by filtering. In those embodiments in which the composition is compounded within an extruder, a fine-mesh screen is typically located at or near the die head of the extruder such that the melt within the extruder must pass through the screen before it leaves the extruder. As the filter becomes plugged with particulates, pressure builds within the extruder and the operational efficiency of the extruder drops until the filter is cleansed or replaced. In those embodiments in which an oligomeric or high molecular weight base, e.g., a oligomeric or high molecular weight HAS, is present in the composition prior to melt filtration, it tends to contribute to the plugging of the screen and diminishing the operational efficiency of the extruder and overall run efficiency of the process.\nInsulation for use in medium voltage power cable applications can typically tolerate more impurities than those for use in high or extra high voltage power cable applications. As such, the screen used to filter the composition before extrusion into pellets can be more coarse, i.e., have a larger openings, than that used for filtering compositions for use in high or extra high voltage power cable applications. As a consequence and all else being equal, the finer (smaller) the screen mesh through which a melt must pass, the more particulate it will trap, the faster it will plug, and the shorter the time interval will be between filter cleaning and/or replacement. This, in turn, affects the operational efficiency of the compounding process.\nOf particular interest to the extension of compositions currently designed for use in medium voltage power cable applications to high and extra high voltage power cable applications is the reduction and/or elimination of particulate contaminants and gels during the compounding of the base polymer with additives and/or fillers and to the extent that such gels are made, their removal by filtering before the composition is fabricated into pellets. Further to this interest is maintaining the relative stability of the pellet against loss of crosslinking efficiency during shipping and/or storage, and the minimizing of water generation during cure."}
{"text": "The present invention relates to a programmable nonvolatile memory, and a technique effective if applied to, for example, an electrically programmable AND type flash memory.\nAn AND type flash memory has been described in a patent document 1 (Japanese Unexamined Patent Publication No. 2004-152977). As one flash memory, there is shown a structure wherein diffusion regions are repeatedly parallel-formed over a semiconductor substrate, auxiliary electrodes are disposed among the respective diffusion regions through an oxide film interposed thereamong to form control transistors, and memory transistors each based on a charge storage region and a control gate are formed on the right and left sides of the auxiliary electrodes. The control gates extend in the direction of diffusion and intersecting the auxiliary electrodes and function as word lines. Further, there is shown, as another structure, a structure wherein other control transistors using auxiliary electrodes in place of the diffused layers are adopted. When the control transistors are turned on, inversion layers are formed in their channel regions and function as wirings. Since the diffusion regions may not be repeatedly disposed in parallel in the latter structure, the structure is excellent in terms of a further reduction in chip area.\nUpon reading for each memory of the structure, the memory transistor for reading is made conductive to its right and left diffusion regions and inversion layers. At this time, memory information is determined according to whether a change in current occurs in the diffusion region according to the threshold voltage of the memory transistor. Upon writing for the memory of the structure, the memory transistor for writing is made conductive to its right and left diffusion regions and inversion layers to allow a write current to flow from the diffusion regions to the inversion layers. At this time, electric field concentration occurs between the corresponding inversion layer and a channel of the memory transistor by reducing the conductance of the control transistor adjacent to the memory transistor for writing, whereby hot electrons generated by the field concentration are injected into the corresponding charge storage region. This write system is referred to as “non cell-through write system”."}
{"text": "The invention relates generally to encoding and decoding of video information, especially such a used for the transmission of video information, and more particularly to systems and methods for video conferencing.\nThere are many systems and techniques for transmitting video information. The most effective conventional techniques involve special transmission and reception systems and require dedicated communication links to encode, transmit, receive and decode video information. The encoding, transmission, and decoding operations are generally resource intensive in terms of the processing (e.g., memory, CPU speed) and transmission requirements (e.g., communication link bandwidth) necessary to provide an adequate video presentation. Such special systems also are generally expensive to own and operate and therefore are not available to an average consumer.\nConventional methods for sending video data are shown and described in U.S. Pat. Nos. 5,740,278, entitled FACSIMILE-BASED VIDEO COMPRESSION METHOD AND SYSTEM, and 5,973,626, entitled BYTE-BASED PREFIX ENCODING, each of which are incorporated by reference herein by their entirety and form a part of this specification.\nAs discussed in U.S. Pat. No. 5,740,278, and in other conventional systems such as MPEG, video data is encoded in either what is referred to herein as “intra” frames (key frames) that hold a single frame image (which only refers to itself) or “inter” frames (delta or difference frames) that refer to another frame (either an intra frame or another inter frame). There are problems with such methods in that:\n(1) network traffic associated with transmission of the two different types of frames is bursty, sometimes requiring a high available bandwidth and other times requiring only a lower available bandwidth due to the extraordinary amount of data transmitted in intra frames vs. inter frames (data in intra frames represents an entire image and are generally larger than inter frames which include more repetitive data and therefore can be compressed more easily, i.e. data rate and required bandwidth vary widely);\n(2) loss of a single intra frame causes loss of quality, and inter frames subsequent to such a loss have no meaning, as they are only meaningful in reference to an intra frame to which they refer or another inter frame (which can also be lost) to which they refer;\n(3) if an intra frame is lost, recovery is very inefficient because the encoder must transmit a new intra frame to regain synchronization, as the entire frame needs to be transmitted, even to replace only a lost portion of the frame; and\n(4) the encoding technique produces an image which is initially very good (intra frame), but the quality of the image slowly degrades over time until transmission of next intra frame.\nIn conventional fax-based video encoding, first a difference frame is generated representing the difference in intensity of each pixel location in two succeeding image frames. On a line-by-line basis, the difference values are encoded. In order to prevent loss of image data or loss of image synchronization (that is, where each line starts and ends), each line of the transmitted video signal has an end-of-line code appended to the end of each line. In the case of conventional fax systems, this end-of-line code is used to protect against bit errors. Consequently, the end-of-line code is large in order to provide sufficient redundancy to correct random bit errors to the degree desired.\nWith the advent of public networks such as the Internet, many commercial products are available to the average consumer for transmitting video information over public networks. These systems may be, for example, coupled with a personal computer for use over the Internet or other public or private communication network. For example, a video conferencing or video distribution system may be configured to transmit video information over the Internet among a group of PCs. Due to the intense resource requirements necessary for transmitting such information, and the limited and/or unreliable resources available on public networks, however, performance of such systems falls short of expectations, and such systems are rendered less usable than more expensive specialized systems. What is needed, therefore, are improved methods for communicating video information."}
{"text": "In the related art, as a method for detecting a normal-line direction of a curved surface such as in a processing surface of a three-dimensional work piece, a method disclosed in PTL 1 is known.\nIn the method disclosed in PTL 1, distance measurement means for measuring at least three positions in the vicinity of a reference point on a curved surface of an object is provided, a virtual plane including each measurement position is set based on each measurement position, and a normal-line direction of the virtual plane is detected.\nIn the method disclosed in PTL 1, when there is unevenness such as a protrusion or a step on a measuring object, it is not possible to detect a correct normal-line. In addition, in PTL 1, a plane approximation method in a case where four or more measurement positions are used is not specifically disclosed.\nMeanwhile, PTL 2 discloses a method in which distance sensors are arranged so as to be deviated by an eccentric amount R in a radial direction from the rotation (processing) center of a measurement axis, the distance sensors are swiveled about the measurement axis, information related to the distance to a continuous measured surface is acquired, and a normal-line is detected by removing a step from the distance information."}
{"text": "The semiconductor integrated circuit (IC) industry has experienced rapid growth. In the course of IC evolution, functional density (i.e., the number of interconnected devices per chip area) has generally increased while geometry size (i.e., the smallest component (or line) that can be created using a fabrication process) has decreased. This scaling down process generally provides benefits by increasing production efficiency and lowering associated costs. Such scaling down has also increased the complexity of processing and manufacturing ICs, and, for these advances to be realized, similar developments in IC manufacturing are needed.\nFor example, mask overlay has become increasingly important as device size shrinks. ICs are typically assembled by layering features on a semiconductor wafer using a set of photolithographic masks. Each mask in the set has a pattern formed by transmissive or reflective regions. During a photolithographic exposure, radiation such as ultraviolet light passes through or reflects off the mask before striking a photoresist coating on the wafer. The mask transfers the pattern onto the photoresist, which is then selectively removed to reveal the pattern. The wafer then undergoes processing steps that take advantage of the shape of the remaining photoresist to create circuit features on the wafer. When the processing steps are complete, photoresist is reapplied and wafer is exposed using the next mask. In this way, the features are layered to produce the final circuit.\nRegardless of whether a mask is error-free, if all or part of the mask is not aligned properly, the resulting features may not align correctly with those on adjoining layers. This misalignment can result in reduced device performance or complete device failure. Conventional overlay metrology tools are used to check mask alignment but have not been satisfactory in all regards. The tools are expensive, slow, and are limited in their ability to detect variations within the die area."}
{"text": "The present invention relates to a circular comb for combing machines with a segment-shaped basic element on which the individual needles (teeth) are located parallel to each other and are connected by pressure strips fastened to the basic element under pressure.\nIn modern combing technology, the teeth (needles) are no longer fastened to a barrette (or bar), but to a wire or strip-like needle carrier to which they are soldered, welded or glued. For fastening such needle strips to a circular comb in a known embodiment (e.g. German Patent DT-OS No. 2,002,020), wedge-shaped grooves are located on the outside circumference of the basic element at predetermined intervals. Needle strips are inserted into these grooves and are pressed by means of a wedge-shaped clamping strip against the webs of the basic element remaining between the grooves. The necessary pressure is achieved by a number of screws distributed in the lengthwise direction over the clamping strip. These screws are threaded into the basic element. Since the combs are relatively close together and the needle tips in the peripheral direction are only 8 mm apart at the outside circumference of the basic element, and since the clamping strips are tapering conically inward, only very small screws can be used for fasteners. The same applies with respect to the pressure screws for detaching the clamping strips from the basic element; additional threads for these must be located in the clamping strip. When replacing the needle strips, a large number of small screws must be unscrewed and tightened again; this requires a considerable expenditure of time.\nIt is also known in the art how to screw fasten the individual needle strips one after the other to the basic element of a circular comb. The first needle strip is placed on the outside of the basic element and fastened by means of a wedge-shaped strip through which the screws pass to engage the threads in the basic element. The needle strip is pressed between the basic element and the wedge-shaped strip. This procedure is repeated till the last needle strip of the segment is reached; it is followed by a final segment which is designed so that the circular comb can be mounted on the machine shaft. The manufacture of such a circular comb is expensive. Even greater is the disadvantage that the individual needle strips cannot be exchanged one for one. In the extreme case, all preceding needle strips of such a circular comb must be removed before the last needle strip can be detached and replaced. This results in extremely cumbersome shutdowns.\nIt is, therefore, an object of the present invention to simplify the fastening of the needle strip on the circular comb of a combing machine and to provide individual interchangeability of the individual needle strips, regardless of whether it is the first or the last or any other needle strip.\nAnother object of the present invention is to provide a circular comb arrangement which may be economically fabricated and maintained in service.\nA further object of the present invention is to provide a circular comb arrangement, as described, which has a substantially long operating life."}
{"text": "Universal Serial Bus (USB) technology allows numerous peripheral devices to be connected to computing devices in a plug-and-play fashion. Such devices include, for example, keyboards, speakers, cameras, joysticks, mice, hard drives, flash drives, DVD drives, and various transceivers. Current peripheral devices are designed and implemented as defined by the Universal Serial Bus 2.0 Specifications, Revision 2.0, which is herein incorporated by reference in its entirety. Users now expect a high level of performance from USB devices. These peripheral devices require ever-increasing bus bandwidth. Therefore, USB technology is evolving from USB 2.0 “High-Speed” to USB 3.0 “SuperSpeed”.\nThe Universal Serial Bus 3.0 Specifications, Revision 1.0, which is also herein incorporated by reference in its entirety, define a number of criteria to be met in order to comply with the USB 3.0 Specifications. USB 3.0 improves on USB 2.0 by improving power management while leveraging existing USB infrastructure. USB 3.0 is a physical SuperSpeed bus combined in parallel with a physical USB 2.0 bus. It has similar architectural components as USB 2.0, including USB 3.0 interconnect, USB 3.0 devices, and USB 3.0 host. The USB interconnect is the manner in which USB 3.0 and USB 2.0 devices connect to and communicate with the USB 3.0 host. The USB 3.0 interconnect inherits core architectural elements from USB 2.0, although several are modified to accommodate the dual bus architecture. Modifications in USB 3.0 include eight primary conductors: three twisted signal pairs for USB data paths and a power pair. One of the twisted signal pairs accommodates for USB 2.0 data path, while two of the twisted signal pairs are used to provide USB 3.0 data paths, one for the transmit path and one for the receive path. In all, USB 3.0 inherits the Vbus, D+, D−, and GND wires from USB 2.0, and incorporates VDD33 conductors to accommodate for SuperSpeed interfaces. USB 3.0 accommodates forwards and backwards-compatibility with existing USB 2.0 peripherals at a lower speed using a Type-A connector.\nWhile USB technology evolves towards the USB 3.0 standard, many current computing devices and peripherals only support USB 2.0. One such peripheral includes a transceiver. 40-nm FPGA's and ASCI's with transceivers have higher integration than prior nodes, including the 65-nm and 45-nm nodes. Another performance benefit of the 40-nm process includes shorter minimum transistor gate lengths than the 45-nm process. Further, power consumption is reduced in the 40-nm node, as smaller process geometries reduce parasitic capacitances. Therefore a need exists for integrating a USB 2.0 transceiver on the same SOC as a USB 3.0 PHY without incurring excess area or system costs."}
{"text": "Embodiments of the present invention generally relate to Radio Frequency Identification (RFID) applications. More specifically, embodiments of the present invention relate to techniques for real-time and offline location tracking using passive RFID technologies.\nRadio Frequency Identification (RFID) is an automatic identification method which relies on the storing and remotely retrieving of data using devices, such as RFID tags or transponders. RFID tags or transponders are also known as proximity, proxy, or contactless cards, because data from an RFID tag can be retrieved without physical contact. Generally, a device, such as an RFID reader, uses radio waves to remotely retrieve a unique identifier stored using the RFID tag when the RFID tag is within proximity of the RFID reader. RFID tags can be attached to or incorporated into a product, animal, or person for the purpose of identification by the RFID reader. RFID readers can be placed on doorways, in train cars, over freeways, mounted on vehicles, and also can be embodied in mobile handheld devices.\nRFID technologies have been traditionally implemented in different ways by different manufacturers, although global standards are being developed. Thus, computer applications using RFID are also typically hard-coded to specific RFID devices sold by the same manufacture. One problem with this arrangement is that the computer applications are limited to using only the data retrieved from the specific RFID readers.\nIn order to provide automated shipping and receiving, real-time inventory, automated shipping and received, and real-time security, other types of RFID sensor devices, such as environment sensors (e.g., temperature and humidity sensors), location sensors (e.g., Global Positioning System or GPS devices), and notification devices, may be required. For example, one cold chain solution provides an RFID tag embedded with a temperature sensor. Cold chain refers to a temperature-controlled supply chain. An unbroken cold chain is an uninterrupted series of storage and distribution activities which maintain a given temperature range.\nHowever, one problem with embedding sensors with RFID tags is that the increase in cost and complexity associated with each RFID tag. Furthermore, if computer applications are tied directly to specific RFID readers, the only items for which sensor data can be used from those applications are those that can be tagged and directly sensed using the specific RFID readers.\nAccordingly, what is desired are improved methods and apparatus for solving the problems discussed above, while reducing the drawbacks discussed above."}
{"text": "Compression ignition engines cause combustion of a hydrocarbon by injecting the hydrocarbon into compressed air and can be fuelled by diesel fuel, biodiesel fuel, blends of diesel and biodiesel fuels and compressed natural gas. The purpose of the present invention is different from the invention claimed in UK patent application no. 1003244.9 filed on 26 Feb. 2010 entitled “Filter”. The purpose of the invention in that patent application is a filter for particulate matter in exhaust gas of a positive ignition engine.\nAmbient PM is divided by most authors into the following categories based on their aerodynamic diameter (the aerodynamic diameter is defined as the diameter of a 1 g/cm3 density sphere of the same settling velocity in air as the measured particle):                (i) PM-10—particles of an aerodynamic diameter of less than 10 μm;        (ii) Fine particles of diameters below 2.5 μm (PM-2.5);        (iii) Ultrafine particles of diameters below 0.1 μm (or 100 nm); and        (iv) Nanoparticles, characterised by diameters of less than 50 nm.        \nSince the mid-1990's, particle size distributions of particulates exhausted from internal combustion engines have received increasing attention due to possible adverse health effects of fine and ultrafine particles. Concentrations of PM-10 particulates in ambient air are regulated by law in the USA. A new, additional ambient air quality standard for PM-2.5 was introduced in the USA in 1997 as a result of health studies that indicated a strong correlation between human mortality and the concentration of fine particles below 2.5 μm.\nInterest has now shifted towards nanoparticles generated by diesel and gasoline engines because they are understood to penetrate more deeply into human lungs than particulates of greater size and consequently they are believed to be more harmful than larger particles, extrapolated from the findings of studies into particulates in the 2.5-10.0 μm range.\nSize distributions of diesel particulates have a well-established bimodal character that corresponds to the particle nucleation and agglomeration mechanisms, with the corresponding particle types referred to as the nuclei mode and the accumulation mode respectively (see FIG. 1). As can be seen from FIG. 1, in the nuclei mode, diesel PM is composed of numerous small particles holding very little mass. Nearly all diesel particulates have sizes of significantly less than 1 μm, i.e. they comprise a mixture of fine, i.e. falling under the 1997 US law, ultrafine and nanoparticles.\nNuclei mode particles are believed to be composed mostly of volatile condensates (hydrocarbons, sulfuric acid, nitric acid etc) and contain little solid material, such as ash and carbon. Accumulation mode particles are understood to comprise solids (carbon, metallic ash etc.) intermixed with condensates and adsorbed material (heavy hydrocarbons, sulfur species, nitrogen oxide derivatives etc.). Coarse mode particles are not believed to be generated in the diesel combustion process and may be formed through mechanisms such as deposition and subsequent re-entrainment of particulate material from the walls of an engine cylinder, exhaust system, or the particulate sampling system. The relationship between these modes is shown in FIG. 1.\nThe composition of nucleating particles may change with engine operating conditions, environmental condition (particularly temperature and humidity), dilution and sampling system conditions. Laboratory work and theory have shown that most of the nuclei mode formation and growth occur in the low dilution ratio range. In this range, gas to particle conversion of volatile particle precursors, like heavy hydrocarbons and sulfuric acid, leads to simultaneous nucleation and growth of the nuclei mode and adsorption onto existing particles in the accumulation mode. Laboratory tests (see e.g. SAE 980525 and SAE 2001-01-0201) have shown that nuclei mode formation increases strongly with decreasing air dilution temperature but there is conflicting evidence on whether humidity has an influence.\nGenerally, low temperature, low dilution ratios, high humidity and long residence times favour nanoparticles formation and growth. Studies have shown that nanoparticles consist mainly of volatile material like heavy hydrocarbons and sulfuric acid with evidence of solid fraction only at very high loads.\nParticulate collection of diesel particulates in a diesel particulate filter is based on the principle of separating gas-borne particulates from the gas phase using a porous barrier. Diesel filters can be defined as deep-bed filters and/or surface-type filters. In deep-bed filters, the mean pore size of filter media is bigger than the mean diameter of collected particles. The particles are deposited on the media through a combination of depth filtration mechanisms, including diffusional deposition (Brownian motion), inertial deposition (impaction) and flow-line interception (Brownian motion or inertia).\nIn surface-type filters, the pore diameter of the filter media is less than the diameter of the PM, so PM is separated by sieving. Separation is done by a build-up of collected diesel PM itself, which build-up is commonly referred to as “filtration cake” and the process as “cake filtration”.\nIt is understood that diesel particulate filters, such as ceramic wallflow monoliths, may work through a combination of depth and surface filtration: a filtration cake develops at higher soot loads when the depth filtration capacity is saturated and a particulate layer starts covering the filtration surface. Depth filtration is characterized by somewhat lower filtration efficiency and lower pressure drop than the cake filtration.\nSelective catalytic reduction (SCR) of NOx by nitrogenous compounds, such as ammonia or urea, was first developed for treating industrial stationary applications. SCR technology was first used in thermal power plants in Japan in the late 1970s, and has seen widespread application in Europe since the mid-1980s. In the USA, SCR systems were introduced for gas turbines in the 1990s and have been used more recently in coal-fired powerplants. In addition to coal-fired cogeneration plants and gas turbines, SCR applications include plant and refinery heaters and boilers in the chemical processing industry, furnaces, coke ovens, municipal waste plants and incinerators. More recently, NOx reduction systems based on SCR technology are being developed for a number of vehicular (mobile) applications in Europe, Japan, and the USA, e.g. for treating diesel exhaust gas.\nSeveral chemical reactions occur in an NH3 SCR system, all of which represent desirable reactions that reduce NOx to nitrogen. The dominant reaction is represented by reaction (1).4NO+4NH3+O2→4N2+6H2O  (1)\nCompeting, non-selective reactions with oxygen can produce secondary emissions or may unproductively consume ammonia. One such non-selective reaction is the complete oxidation of ammonia, shown in reaction (2).4NH3+5O2→4NO+6H2O  (2)\nAlso, side reactions may lead to undesirable products such as N2O, as represented by reaction (3).4NH3+5NO+3O2→4N2O+6H2O  (3)\nVarious catalysts for promoting NH3-SCR are known including V2O5/WO3/TiO2 and transition metal/zeolites such as Fe/Beta (see U.S. Pat. No. 4,961,917) and transition metal/small pore zeolites (see WO 2008/132452).\nEP 1663458 discloses a SCR filter, wherein the filter is a wallflow monolith and wherein an SCR catalyst composition permeates walls of the wallflow monolith. The specification discloses generally that the walls of the wallflow filter can contain thereon or therein (i.e. not both) one or more catalytic materials. According to the disclosure, “permeate”, when used to describe the dispersion of a catalyst slurry on the wallflow monolith substrate, means the catalyst composition is dispersed throughout the wall of the substrate.\nWO 2008/136232 A1 discloses a honeycomb filter having a cell wall composed of a porous cell wall base material and, provided on its inflow side only or on its inflow and outflow sides, a surface layer and satisfying the following requirements (1) to (5) is used as DPF: (1) the peak pore diameter of the surface layer is identical with or smaller than the average pore diameter of the cell wall base material, and the porosity of the surface layer is larger than that of the cell wall base material; (2) with respect to the surface layer, the peak pore diameter is from 0.3 to less than 20 μm, and the porosity is from 60 to less than 95% (measured by mercury penetration method); (3) the thickness (L1) of the surface layer is from 0.5 to less than 30% of the thickness (L2) of the cell wall; (4) the mass of the surface layer per filtration area is from 0.01 to less than 6 mg/cm2; and (5) with respect to the cell wall base material, the average pore diameter is from 10 to less than 60 μm, and the porosity is from 40 to less than 65%. See also SAE paper 2009-01-0292.\nNOx absorber catalysts (NACs) are known e.g. from U.S. Pat. No. 5,473,887 and are designed to adsorb nitrogen oxides (NOx) from lean exhaust gas (lambda>1) and to desorb the NOx when the oxygen concentration in the exhaust gas is decreased. Desorbed NOx may be reduced to N2 with a suitable reductant, e.g. gasoline fuel, promoted by a catalyst component, such as rhodium, of the NAC itself or located downstream of the NAC. In practice, control of oxygen concentration can be adjusted to a desired redox composition intermittently in response to a calculated remaining NOx adsorption capacity of the NAC, e.g. richer than normal engine running operation (but still lean of stoichiometric or lambda=1 composition), stoichiometric or rich of stoichiometric (lambda<1). The oxygen concentration can be adjusted by a number of means, e.g. throttling, injection of additional hydrocarbon fuel into an engine cylinder such as during the exhaust stroke or injecting hydrocarbon fuel directly into exhaust gas downstream of an engine manifold.\nA typical NAC formulation includes a catalytic oxidation component, such as platinum, a significant quantity, i.e. substantially more than is required for use as a promoter such as a promoter in a TWC, of a NOx-storage component, such as barium, and a reduction catalyst, e.g. rhodium. One mechanism commonly given for NOx-storage from a lean exhaust gas for this formulation is:NO+½O2→NO2  (4); andBaO+NO2+½O2→Ba(NO3)2  (5),wherein in reaction (4), the nitric oxide reacts with oxygen on active oxidation sites on the platinum to form NO2. Reaction (5) involves adsorption of the NO2 by the storage material in the form of an inorganic nitrate.\nAt lower oxygen concentrations and/or at elevated temperatures, the nitrate species become thermodynamically unstable and decompose, producing NO or NO2 according to reaction (6) below. In the presence of a suitable reductant, these nitrogen oxides are subsequently reduced by carbon monoxide, hydrogen and hydrocarbons to N2, which can take place over the reduction catalyst (see reaction (5)).Ba(NO3)2→BaO+2NO+ 3/2O2 or Ba(NO3)2→BaO+2NO2+½O2  (6); andNO+CO→½N2+CO2  (7);(Other reactions include Ba(NO3)2+8H2→BaO+2NH3+5H2O followed by NH3+NOx→N2+yH2O or 2NH3+2O2+CO→N2+3H2O+CO2 etc.).\nIn the reactions of (4)-(7) above, the reactive barium species is given as the oxide. However, it is understood that in the presence of air most of the barium is in the form of the carbonate or possibly the hydroxide. The skilled person can adapt the above reaction schemes accordingly for species of barium other than the oxide and sequence of catalytic coatings in the exhaust stream.\nIn Europe, since the year 2000 (Euro 3 emission standard) emissions are tested over the New European Driving Cycle (NEDC). This consists of four repeats of the previous ECE 15 driving cycle plus one Extra Urban Driving Cycle (EUDC) with no 40 second warm-up period before beginning emission sampling. This modified cold start test is also referred to as the “MVEG-B” drive cycle. All emissions are expressed in g/km.\nThe Euro 5/6 implementing legislation introduces a new PM mass emission measurement method developed by the UN/ECE Particulate Measurement Programme (PMP) which adjusts the PM mass emission limits to account for differences in results using old and the new methods. The Euro 5/6 legislation also introduces a particle number emission limit (PMP method), in addition to the mass-based limits.\nEmission legislation in Europe from 1 Sep. 2014 (Euro 6) requires control of the number of particles emitted from both diesel and gasoline passenger cars. For diesel EU light duty vehicles the allowable limits are: 500 mg/km carbon monoxide; 80 mg/km nitrogen oxides (NOx); 170 mg/km total hydrocarbons+NOx; 4.5 g/km particulate matter (PM); and particulate number standard of 6.0×1011 per km. The present specification is based on the assumption that this number will be adopted in due course.\nA difficulty in coating a filter with a catalyst composition is to balance a desired catalytic activity, which generally increases with washcoat loading, with the backpressure that is caused by the filter in use (increased washcoat loading generally increases backpressure) and filtration efficiency (backpressure can be reduced by adopting wider mean pore size and higher porosity substrates at the expense of filtration efficiency)."}
{"text": "The present invention relates to communication of data and in particular to modulation of data transmitted between two entities of a communication system.\nA communication system is a facility which facilitates communication between two or more entities such as communication devices, network entities and other nodes. A communication system may be provided by one more interconnect networks. It is noted that although a communication system typically comprises at least one communication network, for example a fixed line network or a wireless or mobile network, in its simplest form a communication system is provided by two entities communicating with each other. The communication may comprise, for example, communication of data for carrying communications such as voice, electronic mail (email), text message, multimedia and so on.\nA user may communicate by an appropriate communication device. An appropriate access system allows the communication device to access to a communication system. An access to the communications system may be provided by a fixed line or wireless communication interface, or a combination of these. Examples of wireless access systems include cellular access networks, various wireless local area networks (WLANs), wireless personal area networks (WPANs), satellite based communication systems and various combinations of these.\nA communication system typically operates in accordance with a standard and/or a set of specifications and protocols which set out what the various elements of the system are permitted to do and how that should be achieved. For example, it is typically defined if the user, or more precisely user device, is provided with a circuit switched bearer or a packet switched bearer, or both. Also, the manner in which communication and various aspects thereof should be implemented between the user device and the various elements of the communication and their functions and responsibilities are typically defined by a predefined communication protocol.\nIn a radio or wireless system an entity in the form of a base station provides a node for communication with user communication devices, often referred to as mobile stations. Communications in the direction from the base station to the user device is seen as occurring on a “downlink” (DL). Communications in the direction from the user device to the base station is then seen as occurring on an “uplink” (UL). It is noted that in certain systems a base station is called ‘Node B’.\nSignalling between various entities may be divided between signalling of control data and actual data. The latter refers to the data content the users wish to communicate. Control signalling, in turn, associates to transfer of information that is not related as such to the data content the users may wish to be transferred. In the following these two forms of signalling are separated by referring to control signalling and data signalling, where appropriate.\nTo ensure proper operation of the system, the control signalling typically has higher quality requirements than the data signalling. On the other hand, the amount of information conveyed by control signalling such as by acknowledgement signalling is typically only one or two bits. This is less than what can be carried by data modulated symbols, for example Quadrature amplitude modulated (QAM) symbols. For example 16QAM carries 4 bits and 64 QAM carries 6 bits.\nAn example of control signalling is the transfer of positive and negative acknowledgement information signalling, often referred to as ACK/NACK signalling. The acknowledgement signalling is used to provide feedback concerning previous transmissions, for example if a previous data transmission such as a data packet is properly received.\nDespite the advantages in signalling technologies, there is still need to optimize the performance of communications between two devices, for example though reduction of errors in control signalling. Use of a single modulation method for control signalling and data signalling might be desired in various applications.\nThe herein described embodiments aim to address one or several of the above mentioned shortcomings and/or desires."}
{"text": "The making of dental impressions is a precise art because of the necessity of forming a mold in which an accurate model of dental anatomy may be made. Heretofore, it has been customary for many years to mix the material in which the impression is to be made from certain ingredients and then disposing the same in a conventional impression tray of various types. It is obvious that if the material in which the impression is to be made is relatively readily flowable, it is capable of entering fissures and interstices, whereby, when the mold material is introduced into the molded cavity of the impression, it will be capable of reproducing the minute as well as the major features of the dental anatomy. Obviously; however, the more viscous the impression material, the less likelihood there is of small details of the anatomy being reproduced in the mold of which the pattern is to be made.\nExamples of previous techniques in the use of dental impression material in an impression tray or otherwise, heretofore, are found in the prior art U.S. Pat. No. 3,390,458 to Lytton dated July 2, 1968, represents a special type of device for making dental impressions particularly capable of pressing the surrounding gum away from a tooth around the gingival, whereby the pattern molded from the impression will extend below the gum line.\nAnother example of the use of a conventional tray is the subject matter of U.S. Pat. No. 3,552,601 to Jahn dated May 13, 1975. In this patent, there is provided a spacer hood which is stretched over the impression material and the impression tray is inserted into the mouth of a patient and the patient is asked to close his mouth in biting on the impression material. To form a more precise impression, the spacer hood then is removed, a secondary impression material is applied and the impression tray again is introduced into the mouth of the patient so that a second impression may be taken.\nA more recent development in the production of loaded dental impression trays which are pre-filled with impression material comprises the subject matter of U.S. Pat. No. 4,553,936 to Wang dated Nov. 19, 1985, and assigned to the assignee of the present invention. In this disclosure a transparent impression tray is provided which is filled with a light-curable impression material. When thus filled, the impression material is covered with a light-opaque covering which extends across the top and ends of the tray and the exterior surfaces of the tray are covered with light-opaque material such as metal foil; the foil serving an additional advantage of reflecting actinic light which is applied to the transparent tray and thus causes the light to permeate the entire mass of the impression material when such artificial light is applied thereto.\nThe prior art methods of using dental impression materials are cumbersome to use in the dental operatory or laboratory because many of the prior art materials are cured using two-component self-cured systems, and must therefore be prepared in the operatory or laboratory immediately prior to use. It is difficult for the practitioner to keep air voids out of the material when it is mixed; and after the preparation has begun, he has a limited time, usually about 5 minutes, in which to use the material before it self cures or begins to cure. Consequently, batches of dental impression materials made in the operatory and laboratory sometimes have to be discarded and the procedures repeated\nAlso, it is sometimes difficult for the practitioner to judge the correct amount of material to be used to obtain an impression, and excess material has a tendency to escape from the tray and become loose in the mouth, and said loose material may cause the patient discomfort and trigger the gagging reflex. When flowable impression materials are used, the tendency of the material to flow out of a dental tray may cause similar problems.\nAccordingly, there is a need in the art for means and a method for obtaining dental impressions whereby the practitioners time in the operatory is reduced, there is less mess involved, the amounts of impression material to be used are premeasured, the dental impression material to be used is free of air voids, where means are provided to prevent dental impression material from flowing out of a tray, and provides for increased patient comfort\nIt has been found that the above objectives can be obtained by using a light activated, prepackaged, premeasured impression material"}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a conductive polymer electrolytic capacitor, and more specifically, it relates to a method of manufacturing a conductive polymer electrolytic capacitor comprising a solid electrolyte containing a conductive polymer and an ionic liquid.\n2. Description of the Background Art\nAn electrolytic capacitor employing a conductive polymer for an electrolyte and a cathode conductive layer is known in general. In relation to such a conductive polymer electrolytic capacitor, it is known that an electrolytic capacitor exhibiting a low leakage current and having high heat resistance and high moisture resistance can be manufactured by employing a dopant hardly damaging a dielectric layer consisting of an anodized film or also using a solid organic onium salt having repairability for a valve metal (refer to Japanese Patent Laying-Open No. 2003-22938, for example).\nHowever, the conductive polymer essentially has no anodic oxidizability, whereby improvement in the withstand voltage characteristic of the conductive polymer electrolytic capacitor is limited. Japanese Patent Laying-Open No. 3-96210 (1991) describes a method of improving the withstand voltage by aging an anodic body having a solid electrolytic layer thereby performing re-repair/forming, as means for solving this problem.\nHowever, this document, describing a method of aging the anodic body by applying a constant voltage not more than half a forming voltage, discloses that the yield is remarkably deteriorated if the forming voltage is increased. Thus, there is no technique disclosing that at least 50% of a forming voltage for a valve metal can be extracted as a normal voltage in relation to a conductive polymer electrolytic capacitor.\nFurther, it is known that dissociation between a forming voltage and a withstand voltage is increased when the forming voltage is increased (refer to Electrolytic Condenser Review, Vol. 53 (1), 95 (2002)). While the forming voltage and the withstand voltage are equivalent to each other up to about 30 V (however, the actual working voltage is set to about 15 V in consideration of safety), the withstand voltage is remarkably reduced at a forming voltage exceeding 30 V. The withstand voltage is 50 V and the practical voltage is not more than 24 V whether forming is performed at 100 V or 300 V. When the forming voltage is increased, further, the capacitance is unpractically reduced. In general, therefore, it is extremely difficult to manufacture an electrolytic capacitor having an actual working voltage of at least 24 V in relation to a conductive polymer electrolytic capacitor.\nIn order to solve this problem, it is attempted to provide an insulating layer referred to as a buffer layer on a dielectric film. If such a layer is provided, however, equivalent series resistance (ESR) or a tan δ characteristic is deteriorated, to damage the high performance of the conductive polymer electrolytic capacitor."}
{"text": "1. Field of the Invention\nThe present invention relates to an arc tube, and a method of fabricating the arc tube for a discharge lamp in which a cylindrical shroud glass is welded on and integrated with an arc tube main body having a sealed glass bulb that serves as a discharge portion and is formed at a portion of the main body along the longitudinal direction thereof.\n2. Description of the Related Art\nConventionally, as shown in FIG. 12, an arc tube is configured in a manner that a cylindrical shroud glass 8 for shielding ultraviolet rays is welded on and integrated with an arc tube main body 1 having a sealed glass bulb 2 which serves as a discharge portion and is formed at a portion of the main body along the longitudinal direction thereof, and the sealed glass bulb 2 is covered by the shroud glass 8. Reference numerals 8a, 8b depict welding portions of the shroud glass 8.\nElectrodes a, a are provided in an opposite manner within the sealed glass bulb 2 sandwiched between the pinch seal portions 3a, 3b, and lead wires c, c coupled to molybdenum foils b, b are drawn from the pinch seal portions 3a, 3b at both ends of the glass bulb, respectively. Cylindrical portions 4a, 4b as non-pinch seal portions are formed at the front and rear portions of the pinch seal portions 3a, 3b so as to be extracted therefrom, respectively.\nThe shroud glass 8 cuts ultraviolet rays in a wavelength range that may be harmful to the human body among light emitted from the sealed glass bulb 2.\nA sealed space 7, formed by the shroud glass 8 and surrounding the arc tube main body 1, suppresses devitrification generated at the arc tube. That is, since the lamp room in which the arc tube is disposed communicates with the outside of the lamp room through an air hole performing breathing operation, and the atmosphere within the lamp room contains a lot of moisture, the moisture causes the devitrification generated at the arc tube. Therefore, the arc tube main body 1 is covered by the sealed space 7 so that the arc tube main body 1 does not contact the atmosphere containing the moisture, thereby suppressing the generation of the devitrification.\nIn order to fabricate the arc tube shown in FIG. 12, first, the rod-shaped arc tube main body 1 having cylindrical portions 4a, 4b formed at both ends thereof is fabricated. Thereafter, the arc tube main body 1 is inserted within a shroud glass tube 9, then the front and rear end sides of the shroud glass tube 9 are heated to be molten and softened. After that, the softened portions are deformed by using forming rolls in a direction for reducing the diameter of the shroud glass tube (a direction shown by the arrows in FIG. 12) and pressed against the pinch seal portions 3a, 3b of the arc tube main body 1 at the inside of the glass tube and welded at the pinch seal portions. Then, the shroud glass tube 9 is cut at predetermined portions as necessary.\nHowever, according to the conventional arc tube described above, there arises a problem that the devitrification phenomenon occurs despite the fact that the shroud glass 8 (the shroud glass tube 9) is welded to the arc tube main body 1 to form the sealed space 7.\nThe inventors of the present invention inspected the cause of the occurrence of the devitrification phenomenon and determined that the cause resides in the sectional shape of the arc tube main body (the pinch seal portions 3a, 3b) for welding the shroud glass 8 thereon. That is, although the cross section of the shroud glass tube 9 is cylindrical, the cross section of the pinch seal portion 3a (3b) is rectangular as shown in FIG. 13(a) since it is typically pinched by a pincher. Thus, in the welding process of the shroud glass, as shown by a phantom line in FIG. 13(a), when the shroud glass tube 9, which is molten and softened and pressed in the direction for reducing the diameter thereof along its radius direction, contacts closely to the surface (flat surface) of the pinch seal portion 3a, an opening S extending in the axial direction along the contact surface is formed (see FIG. 13(b)). As a result, the atmosphere (moisture) within the lamp room enters into the sealed space 7 around the arc tube main body 1 from the opening S formed at the welded portion, thereby causing devitrification.\nThe inventors of the present invention have determined that an opening is not formed at the contact surface between the arc tube main body 1 and the shroud glass 8 when the welded portion of the shroud glass at the arc tube main body 1 is formed in a circular shape in its cross section.\nAccordingly, an object of the present invention is to provide an arc tube for a discharge lamp and a method for fabricating the arc tube in which the welded portion of the shroud at an arc tube main body is formed in a circular shape in its cross section thereby preventing the forming of an opening at the contact surface between the arc tube main body and the shroud glass.\nIn order to attain the aforesaid object, an arc tube for a discharge lamp according to the present invention is arranged in a manner such that the arc tube includes an arc tube main body at which a sealed bulb, for example, a glass bulb, serving as a discharge portion sandwiched by front and rear pinch seal portions is formed at a portion of a tube along a longitudinal direction thereof, and a cylindrical shroud which is welded on and integrated with the arc tube main body so as to cover the sealed bulb to form an airtight sealed space around the arc tube main body,\nthe front and rear end portions of the shroud are welded on shroud weld portions with circular cross sections provided at front and rear end sides of the arc tube main body, respectively. In the embodiments described herein, the shroud may be formed of glass, and the arc tube body may be formed from a glass tube. of course, other materials known to those skilled in the art may be substituted without departing from the scope of the present invention.\nA method of fabricating an arc tube for a discharge lamp according to the present invention includes an arc tube main body fabricating process for fabricating an arc tube main body at which a sealed glass bulb serving as a discharge portion sandwiched by front and rear pinch seal portions is formed at a portion of a glass tube along a longitudinal direction thereof, and a shroud glass welding process for welding and integrating a cylindrical shroud glass on and with the arc tube main body so as to cover the sealed glass bulb, wherein\nin the arc tube main body fabricating process, shroud glass welded portions with circular cross sections are formed on front and rear end sides of the arc tube main body, respectively, and\nin the shroud glass welding process, the arc tube main body is inserted into the shroud glass tube, predetermined positions of the shroud glass having been heated, molten and softened are modified in a direction of reducing diameters thereof, and the predetermined positions are welded on the shroud glass welded portions on the front and rear end sides of the arc tube main body, respectively.\nAt the time of welding the shroud glass to the arc tube main body, the predetermined positions of the shroud glass tube having been heated, molten and softened are modified inside so as to reduce their diameter in a radial direction. As shown in FIG. 9 (a diagram showing a state where the shroud glass is shrink-sealed) in an embodiment of the invention, each of the outer peripheral surfaces of the shroud glass welded portion of the arc tube main body (a shrink seal portion 15a and a cylindrical portion 14a) and the inner peripheral surface of the molten diameter-reduced area of the shroud glass tube 20 may have a circular shape almost matching to each other. Thus, the inner peripheral surface of the molten diameter-reduced portion of the shroud glass tube 20 is molten and welded on the outer peripheral surface of the shroud glass welded portion of the arc tube main body (the shrink seal portion 15a and the cylindrical portion 14a) uniformly along the peripheral direction thereof without causing any space therebetween, so that such a space for releasing the sealed space around the arc tube main body to the atmosphere is not formed at the welding portion between the arc tube main body and the shroud glass tube.\nIn particular, when inactive gas, adjusted to be a negative pressure such that a pressure becomes about 1 atm. upon lighting and heating the arc tube, is supplied within the airtight sealed space formed around the arc tube main body, the arc tube main body is prevented from contacting moisture in the atmosphere.\nFurther, in a method of fabricating an arc tube according to the present invention, in the arc tube main body fabricating process, a cylindrical non-pinch seal portion is formed in an extended manner at a backward portion of the pinch seal portion on the rear end side of the arc tube main body, and a shrink seal portion is formed adjacent to a forward portion of the pinch seal portion on the front end side of the arc tube main body, and\nin the shroud glass welding process, the rear end side of the shroud glass tube is welded on the cylindrical non-pinch seal portion on the rear end side of the arc tube main body, and the front end side of the shroud glass tube is welded on the shrink seal portion on the front end side of the arc tube main body.\nAt the rear end portion of the shroud glass, the circular inner peripheral surface on the rear end side of the shroud glass tube which is molten, softened and modified in a direction of reducing the diameter thereof matches almost with the outer peripheral surface of the cylindrical non-pinch seal portion on the arc tube main body side, and so the inner peripheral surface of the molten shroud glass is molten and adhered to the outer peripheral surface of the non-pinch seal portion uniformly along the peripheral direction thereof without causing any space therebetween.\nIn contrast, at the front end portion of the shroud glass, since the outer peripheral surface of the shrink seal portion has a circular shape, the circular inner peripheral surface on the front end side of the shroud glass tube which is molten, softened and modified in a direction of reducing the diameter thereof adheres to the circular outer peripheral surface of the shrink seal portion on the arc tube main body side uniformly along the peripheral direction thereof without causing any space therebetween.\nIncidentally, the welded portion with a circular cross section on the front end side of the arc tube main body may be, for example, a cylindrical portion (see FIG. 11) serving as a non-pinch seal portion extending forward of the front end side pinch seal portion; a pinch seal portion with a circular cross section provided adjacent to the forward portion of the front end side pinch seal portion with a rectangular cross section; a shrink seal portion provided adjacent to the forward portion of the front end side pinch seal portion; the pinch seal portion with the circular cross section and the cylindrical portion (see FIG. 10); or the shrink seal portion and the cylindrical portion (see FIG. 1). The shrink seal portion to which the shroud tube is welded can be formed in the following manner.\nThe arc tube main body may be fabricated in a manner that the predetermined filling material is supplied to the glass bulb of the glass tube which one end side is subjected to the primary pinch sealing, and thereafter the other side of the glass tube is subjected to the secondary pinch sealing. Then, the secondary pinch sealing process is performed in a manner that the seal expected area near the glass bulb is heated, molten and softened while the glass bulb of the glass tube is cooled by using cooling medium. In this respect, prior to the secondary pinch sealing using a pincher, the seal expected area having been heated, molten and softened deforms and shrinks in the diameter reducing direction due to the negative pressure within the glass tube (the negative pressure formed by condensing the filling material such as inactive gas etc. within the glass bulb) and so the shrink seal portion with the circular cross section is formed. In other words, the secondary pinch seal expected area of the glass bulb is entirely shrink-sealed. Then, the glass bulb side of the shrink seal portion is pinch-sealed with the predetermined width (a portion of the shrink seal portion closer to the glass bulb is pinch-sealed so that the shrink seal portion with the predetermined width remains), whereby the shrink seal portion with the circular cross section (shroud glass welded portion) is formed adjacent to the pinch seal portion with the rectangular cross section.\nThe width (length) of the shrink seal portion serving as the shroud glass welded portion may be in a range of L/6 to L/2, where L represents the entire length of the seal portion (that is, the pinch seal portion and the shrink seal portion). The inventors have determined that when the width is equal to or less than L/6, it becomes difficult to weld the shroud glass and a space is generated at the welding surface. In contrast, when the width is equal to or more than L/2, the length of the pinch seal portion becomes shorter, so that the property of the adhesion between the glass layer and the electrode assembly at the seal portion may be degraded and the airtightness of the sealed glass bulb may not be secured.\nFurther, in a method of fabricating an arc tube for a discharge lamp according to the invention,\nin the arc tube main body fabricating process, a cylindrical non-pinch seal portion provided with a circular flange portion on an outer periphery thereof is formed in an extended manner at a backward portion of the pinch seal portion on the rear end side of the arc tube main body, and\nin the shroud glass welding process, the rear end side of the shroud glass tube is welded on the circular flange portion on the rear end side of the arc tube main body.\nThe circular flange portion serving as the shroud glass welded portion is disposed closely to the inside of the rear end portion of the shroud glass tube, and the rear end portion of the shroud glass tube having been heated, molten and soften is molten and welded smoothly on the circular flange portion on the inside.\nAlso, the arc tube main body fabricating process may include a glass bulb forming process for forming a glass bulb at a portion of the glass tube; a primary pinch seal process for inserting an electrode assembly from one end side of the glass tube provided with the glass bulb and pinch-sealing a portion near the glass bulb; a sealing and exhausting process for supplying predetermined filling material such as mercury to the glass bulb, inserting an electrode ashy from the other end side of the glass tube and holding the ashy thereat, supplying inactive gas within the glass bulb and pinch-sealing or tipping off an opening end side of the glass tube to seal within the glass tube; and a secondary pinch seal process for pinch-sealing a portion of the glass tube near the glass bulb,\nthe shroud glass welding process includes a process for welding the rear end side of the shroud glass tube on the rear end side of the arc tube main body, and a process of welding the front end side of the shroud glass tube on the front end side of the arc tube main body,\nin the secondary pinch seal process constituting the arc tube main body fabricating process, a seal expected area near the glass bulb is heated and molten to perform shrink sealing while cooling the glass bulb by using cooling medium, thereafter the glass bulb side of the shrink seal portion is pinch-sealed with a predetermined width to form a shrink seal portion adjacent to the pinch seal portion,\nin the shroud glass tube front end side welding process constituting the shroud glass welding process, a pressure within the shroud glass tube which rear end side being welded on the rear end side of the arc tube main body is kept at a negative pressure, a welding expected area on the front end side of the shroud glass tube is heated, molten and softened, and the front end side of the shroud glass tube is shrink-sealed to the shrink seal portion adjacent to the pinch seal portion.\nIn the secondary pinch seal process of the arc tube main body fabricating process, the seal expected area on the front end side of the glass tube having been heated, molten and softened deforms and shrinks in the diameter reducing direction due to the negative pressure within the glass tube (the negative pressure formed by condensing the filling material such as inactive gas, etc. within the glass bulb) and so the shrink seal portion with the circular cross section is formed. Then, the glass bulb side of the shrink seal portion is pinch-sealed thereby to form the shrink seal portion (shroud glass welded portion) adjacent to the forward portion of the pinch seal portion on the front end side of the arc tube main body.\nIn the shroud glass tube front end side welding process of the shroud glass welding process, the welding expected area of the shroud glass tube having been heated, molten and softened deforms and shrinks in the diameter reducing direction due to the negative pressure within the glass tube and is molten and welded on the shrink seal portion with the circular cross section (the shroud glass welded portion) on the front end side of the arc tube main body.\nIn addition, in the method of fabricating an arc tube,\nin the arc tube main body fabricating process, a cylindrical non-pinch seal portion is formed in an extended manner at a forward portion of the pinch seal portion on the front end side of the arc tube main body, and\nin the shroud glass welding process, the front end side of the shroud glass tube is welded on only the cylindrical non-pinch seal portion on the front end side of the arc tube main body or on a welded portion with a circular cross section including the cylindrical non-pinch seal portion.\nSince the front end side of the shroud glass tube is welded on the cylindrical non-pinch seal portion on the front end side of the arc tube main body, the axial length of the welding surface can be made larger."}
{"text": "1. Field of the Invention\nThe present invention relates to a laminate comprising an acrylic resin layer, a thermoplastic resin layer and a propylene-based resin layer, to a structure comprising the laminate and a polypropylene resin substrate, and to an automotive part and a part of household appliances, the parts comprising the foregoing structure.\n2. Description of the Prior Art\nTo produce a molded resin article of good design property by preforming a decorative or colored film by vacuum forming or the like, inserting the film preformed into a mold for injection molding, thereafter injecting a synthetic resin to allow the decorative or colored film to laminate to a part of the surface of the resulting molded article has recently been proposed as a method for improving the design property of the surface of a molded article without performing painting. For instance, Japanese Patent Kohyo Publication No. 2-503077 and Japanese Patent Laid-Open 11-207896 disclose that a colored molded article is produced by use of a laminate constituted of a transparent layer, a colored layer and a substrate, as a paintless film.\nIncidentally, in the case where a molded article of good design properties is produced by such a method, it is general to use a polyolefin resin such as polypropylene as a resin for injection molding which will form a substrate of the molded article. Further, it is important for a molded article to be excellent in scratch resistance, surface luster property, weather resistance and the like. Furthermore, it is required to produce such a molded article at low cost. Therefore, at present, adopted is a method comprising the steps of preforming a laminate comprising an acrylic resin layer as a transparent layer or a colored layer, a polypropylene resin layer as a substrate layer by thermoforming such as vacuum forming into a predetermined shape, inserting the resulting preformed laminate into a mold for injection molding, and injecting, after the insertion, a polyolefin resin. In the above method, the xe2x80x9csubstrate layerxe2x80x9d indicates a layer in the laminate that will come in contact with an injection molding resin which will become a substrate.\nHowever, it has become clear that such a conventional laminate may form cracks therein during its handling, for example, at the time of its setting for thermoforming or at its release from a mold after its forming. Such cracking will become a serious problem because it will deteriorate the appearance of a structure after the injection of a polyolefin resin.\nThe present invention was made for the elimination of the above problem. The object of the present invention is to provide a laminate of good crack resistance wherein the laminate does not form cracks at the time of its setting during its thermoforming or of its release from a mold after its thermoforming.\nThe present invention provides a laminate comprising at least:\na resin layer (A) containing a propylene-based resin (a);\na resin layer (B) containing a thermoplastic resin (b) with a tensile elongation at break of not less than 100%; and\na resin layer (C) containing an acrylic resin (c), wherein the resin layer (A),\nthe resin layer (B) and the resin layer (C) are disposed in layers in this order.\nThe present invention further provides a structure comprising the foregoing laminate and a substrate containing a polyolefin resin, the substrate being laminated to the resin layer (A) of the laminate. Furthermore, the present invention provides an automotive part and a household appliance part comprised of the foregoing laminate structure."}
{"text": "Piezoelectric ceramic plates are used for a variety of electronic components, such as a piezoelectric actuator, which utilize, as a mechanical driving source, displacement or force generated via a piezoelectric phenomenon. With broadened use of piezoelectric actuators, multi-layer piezoelectric actuators from which larger displacement or larger generated force can be obtained at lower voltages have become increasingly used.\nSince conventional piezoelectric ceramic plates develop significant deformation (variation in shrinkage) after firing, such piezoelectric ceramic plates are subjected to processing, such as cutting and polishing, after firing in order to control the shape and dimensions of the piezoelectric ceramic plates into prescribed ranges (see, for example, Patent Literature 1)."}
{"text": "1. Field of the Invention\nThe present invention relates to an errorproof device, and more particularly to the combination of an errorproof device and a modular socket.\n2. Description of Related Art\nA conventional communication modular socket for connection with a modular plug used in a telephone line or a modem does not have the ability to distinguish whether the plug to be inserted into the socket has the appropriate dimension. For example, the currently available RJ 11 or RJ45 plugs are both used for communication devices and respectively have a dimension different from the other. Because the RJ 11 plug has a smaller dimension than that of the RJ 45 plug, the RJ 11 plug may be erroneously inserted into the modular socket (namely the RJ 45 socket) configured to mate for the RJ 45 plug and thus leads on the RJ 45 modular socket may be damaged.\nTo overcome the shortcomings, the present invention tends to provide an improved modular socket having an errorproof device therein to mitigate the aforementioned problems."}
{"text": "1. Field of the Invention\nThe present invention relates to an underground continuous wall job practice applying into soft base anti-seepage, and meanwhile involves to its special drill.\n2. Description of Prior Art\nIn accordance with the conventional underground continuous walls, most of them are used for anti-seeping in foundation of dams, basements, harbor and sluice gate base or both banks of rivers and so on. The job practice is to excavate a trench in the job location firstly, then to pour concrete to construct a continuous proof-water wall in the pit. In this way, too much earthmoving has to be processed in the construction procedure, and too much the concrete is used, the time limit of project is too long and the production cost is very high. And once encountering with groundwater or complex geological structure, the sidewalls of said trench couldn\"\"t be kept in standing so that the construction difficult is increased. If utilizing steel sheet wall or prefabricated concrete plate pile and so on, the investment cost will double and redouble, the anti-seepage effect is also not good as expected. A rapid construction method of concrete shell pile in soft basement and its special drill disclosed in Chinese Patent No.98113070.4 provides a way of constructing a single pile, but not constructing a continuous wall underground conveniently and economically.\nIt is therefore a main object of the present invention to provide a rapid and economic construction way to build an underground shell-pile continuous wall to have better expected anti-seeping effect, meanwhile to provide a special drill for carrying out this job practice.\nThis object is achieved by a job-practice adapting cluster of shell-piles to build an underground wall. Follows steps carry out said job practice:\n1. Attach a circular pile shoe having cutting faces on the bottom end of the barrel core space of each special-drill building cluster of shell-piles.\n2. Locate the first special drill on the job position vertically standing on said pile shoe, and exert pressure on the vibrating head at the tip end of said special drill to press said special drill sinking into the ground to get the desired depth in the soft soil layer, the expelled soil is driven out along the inside wall of inner draining hole of said special drill.\n3. Locate and connect the second special drill with the first special drill by coupling the male connector of the second special drill to the female connector of the first special drill, and process the second special drill as described in step 2 to sink into the desired depth.\n4. Pour the concrete from the loading hopper of the first special drill, vibrate said special drill as pouring, and withdraw up said special drill simultaneously to depart said pile shoe from the drill trunk until the first special drill is dragged out from the ground completely to build the first shell pile.\n5. Put on a new pile shoe on the used first special drill or prepare the third special drill, repeat the step 3 to locate and connect the new special drill to the second special drill, and sink into the designed depth.\n6. Repeat the step 4 to withdraw the second special drill, meanwhile build the second clustering shell pile with the first one.\n7. In the same way, repeat the circular steps until to build enough many shell-piles to construct an underground continuous wall.\nFor archiving above-described job practice, said special drill is designed to comprise of a pair of centered an inner sleeve and an outer sleeve, a barrel core space formed between the outside wall of said inner sleeve and the inside wall of said outer sleeve, and a circular pile shoe fitted into the ring opening at the bottom side of said barrel core space, and a flange retained on the top ends of said inner and outer sleeves, and a vibrating head attached on the top side of said flange; said vibrating head has a draining hole connecting to the inner cave of said inner sleeve, and a loading hopper set upon the outside wall of said outer sleeve approaching the top end connecting to the barrel core space, and a pair of male connecter and female connector are fixed on the outside wall of said outer sleeve axially and parallelly, and their including angle to the central axial should be not less than 60xc2x0.\nDue to applying this job practice, less earthmoving is processed during the construction procedure, less concrete is needed, it can be used for pouring and building an underground continuous shell-pile wall in the real time at the work site, so that the time limit of project is shortcut, and the cost is reduced, and a better anti-seepage effect is archived.\nSaid special drill provided in the present invention has a simple structure; by rapid connecting and locating male and female connectors fixed on said special drill, conjunction of shell-pile is became into easy. Draining the expelled soil can be processed as sinking the special drill so that the sinking drag turbine is less, and earthmoving is less too."}
{"text": "The invention relates to a valve-train of an internal combustion engine with a camshaft that comprises a carrier shaft as well as a cam element that is locked in rotation on this carrier shaft and that can move between two axial positions and that has at least one cam group of directly adjacent cams with different cam lobes and an axial connecting link constructed as a groove with external guide walls for defining two intersecting connecting link pathways, and with an activation pin that can couple into the axial connecting link for moving the cam element in the direction of both connecting link pathways.\nSuch a valve-train assembly that is used for the variable activation of gas-exchange valves by moveable cam elements and in which a single activation pin is sufficient for each cam element, in order to move the cam element in the direction of both connecting link pathways, is already known from DE 101 48 177 A1, which is considered class-forming. In that publication, two cam elements are disclosed with alternatively constructed axial connecting links, wherein the first axial connecting link has a central guide web for forming inner guide walls for the activation pin and the second axial connecting link consists merely of outer guide walls.\nThe latter construction has the advantage that the production expense for the axial connecting link is significantly lower due to the elimination of the guide web. One significant risk with respect to the functional safety of the valve-train assembly in the case of this construction is that, however, the displacement process of the cam element is completely finished, i.e., without incorrect switching, only when the inertia of the mass in motion of the cam element is sufficient to move it into its other end position after passing through the intersection region of the connecting link pathways without forced guidance of the activation pin, that is, to a certain extent, in free fall. A prerequisite for the sufficient inertia of the mass in motion of the cam element is obviously a minimum rotational speed of the camshaft that is directly dependent on the friction between the cam element and the carrier shaft. Displacement of a cam element with a rotational speed below this minimum rotational speed could have the result that the cam element remains “at a half-way point” and a cam follower acting on the gas-exchange valve is simultaneously acted upon by several cams of the cam group in an uncontrolled manner and simultaneously under high mechanical loading. In addition, in this case there is no longer the ability to move the cam element through action of the activation pin later into one of the end positions, because in this case there is no longer axial allocation between the activation pin and the outer guide walls.\nThis functional risk is indeed significantly smaller in the case of the first construction of the axial connecting link with a central guide web whose inner guide walls cause a further accelerating forced guidance if the rotational speed of the cam element is lower than the activation pin. Nevertheless, there is also the risk here that the activation pin does not pathway into the specified connecting link pathway after passing through the intersection region, but instead collides with the end face of the guide web also under high mechanical loading."}
{"text": "Field of the Invention\nThe invention relates to the field of atomic sensors, and more particularly to a physical unit of a chip-scale nuclear magnetic resonance (NMR) gyroscope.\nDescription of the Related Art\nA typical NMR gyroscope includes a physical unit and an electric unit. The physical unit includes: an optical source, optical elements, an atomic vapor chamber, and photodetectors. The physical system of the chip-scale NMR gyroscope is achieved by a micro electro mechanical system (MEMS) processing. However, MEMS suffers from the following problems: 1) The atomic vapor chamber of the MEMS generally adopts a glass-silicon-glass sandwich structure, but the pump light beam and the probe light beam are difficult to orient orthogonally in the light-atom interaction region. 2) The two light beams are provided by two independent semiconductor lasers, and temperatures, frequencies of the output laser beams, and the powers of the two lasers are separately controlled, thereby resulting in large resource consumption; besides performances of the two light beams change independently, resulting in in difficulties in cooperation of the two light beams under working conditions. 3) Although in some atomic vapor chambers of the MEMS adopting the glass-silicon-glass sandwich structure, the pump light beam and the probe light beam are orthogonal to each other within the light-atom interaction region, it is difficult to couple the light into the atomic vapor chamber."}
{"text": "Tiles are often secured to a substrate such as a floor, wall, countertop, or the like using grout/mortar located underneath the tiles and between the tiles. When tiles are being installed, it is often desirable that adjacent tiles are laid in a manner such that top surfaces of the adjacent tiles are level with each other."}
{"text": "Storage devices, in general, have various methodologies for placing and allocating data onto the designated locations of the applicable storage medium, depending generally on the type of storage medium. As such, most storage media have an “address space” that may be used to associate a physical “storage space” on that media and which address space is presented to consumers of the data storage as the available storage. Depending on the storage media, there may be unique challenges and utilizations for managing the address space for the different types of media.\nFor example, a typical flash memory device may in fact have 1.2 TB of physical storage, but only 800 MB of address space available for use by a data consumer; while described below in greater detail, this is to manage a unique property of flash memory associated with the requirement that prior to updating data in a memory cell that already has data therein, the data must first be deleted (not overwritten) and any deletion happens with reduced granularity than is possible by writes. In other words, a complete erase block (which will be denoted as a “row” of memory cells in some examples) must be first deleted to write new data to any currently occupied/in-use memory cell in that erase block. When such data to be deleted or updated co-exists in an erase block with live data (which therefore cannot be erased), flash devices are presented with a problem which they have evolved in various ways to address.\nIn flash devices, a data address space is used to manage this flash-specific issue, particularly in cases when (1) other memory cells within an erase block have data that must be maintained and (2) an increased number of the available rows of memory have at least one cell with live data. In such cases, data must be reallocated to ensure that there remain available rows of memory cells which can be entirely deleted if necessary. The address space is configured to provide a data address for data stored in memory blocks that are originally associated with memory blocks in a row (i.e. erase block), such row also containing other live and possibly non-related data, with memory blocks in other rows that do not suffer the same limitation (meaning that data can be written thereto without deleting the entire row or that the remaining memory blocks in the row can be deleted without destroying other live data). In order to avoid the scenario in which there are no more rows that can be safely deleted, because all or nearly all rows contain at least one block with live data, the flash device may be configured to consistently re-associate the data in the data address space to preserve available rows; in some cases, the flash device is overprovisioned with memory relative to the address space; in some cases, the flash device is also overprovisioned with processing capability, for example with a dedicated high-speed processor to manage the allocation of the data; in some cases, all or a combination of some of these are implemented.\nWhile the above issue relates to flash memory devices, similar translation layers may exist in other forms of data storage devices in which the physical locations, or addresses or other indicators of the actual location of storage locations within the storage media, may be abstracted by virtual addresses (which may be referred to herein as data addresses). For example, in spinning disk devices, the physical data locations are frequently defragmented or collocated within the same track on a disk, or moved to an outer track (upon which data retrieval may be faster), since data located in close proximity speeds performance in spinning disk media. In order to keep track of the data, a translation layer, such as a register that maintains a log or table keeping track of data addresses and physical location indicators is maintained. The data addresses exist in a data address space and, from the perspective of the physical storage media, are not currently managed in order to impact performance of the device.\nFrom the perspective of the data consumer at any layer or abstraction (like for example, a user, a file system, or the forwarding tables on a network switch), all of this translation, and the underlying physical locations are not visible; the data consumer views the storage device as having a capacity that is equal in size to that of the data address space, and in fact may in cases be indistinguishable from such data address space. From the perspective of a data consumer, a flash device, for example, is a storage resource with a capacity that is the same as the address space.\nFlash memory is becoming a widely used storage media due to its improved performance in some respects: fast access speeds, low-power, non-volatile, and rugged operation. Most flash devices comprise a flash translation layer (FTL) that generally comprises an FTL driver that works in conjunction with an existing operating system (or, in some embedded applications, as the operating system) to make linear flash memory, with erase blocks that are larger than individual write blocks, appear to the system like a single memory resource. It does that by doing a number of things. First, it creates “virtual” small blocks of data, or sectors, out of the flash's large erase blocks. Next, it manages data on the flash so that it appears to be “write in place” when in fact it is being stored in different spots in the flash. Finally, FTL manages the flash so there are clean/erased places to store data.\nFile systems, or other data consumers, such as, for example, operating systems on data consuming devices, such as DOS, typically use drivers that perform input and output in structured pieces called blocks. Block devices may include all disk drives and other mass-storage devices on the computer. FTL emulates a block device. The flash media appears as a contiguous array of storage blocks numbered from zero to one less than the total number of blocks. In the example of DOS interacting with a flash memory device, FTL acts as a translation layer between the native DOS BPB/FAT file system and flash. FTL remaps the data to the physical location at which the data is to be written. This allows the DOS file system to treat flash like any other block storage device and remains ignorant of flash device characteristics. FTL appears to simply take the data from the file system and write it at the specified location (sector). In reality, FTL places the data at a free or erased location on the flash media and notes the real location of the data. It may in some cases also invalidate the block that previously contained the block's data (if any), either in the FTL or in a separate data tracking module elsewhere on, or in data communication with, the FTL or the computing processor resources that is running the FTL. If the file system asks for previously written data, it requests the data at the specified data address and the FTL finds and reads back the proper data from the actual physical storage location. Flash media allows only two states: erased and non-erased. In the erase state, a byte may be either all ones (0xFF) or all zeroes (0x00) depending on the flash device. A given bit of data may only be written when the media is in an erase state. After it is written to, the bit is considered dirty and unusable. In order to return the bit to its erase state, a significantly larger block of flash called an Erase Zone (also known as an erase block) must be erased. Flash technology does not allow the toggling of individual bits or bytes from a non-erased state back to an erased state. FTL shields the file system from these details and remaps the data passed to it by writing to unused data areas in the flash media. This presents the illusion to DOS, or other data consumer, that a data block is simply overwritten when it is modified. In fact, the amended data has been written somewhere else on the media. FTL may, in some cases, also take care of reclaiming the discarded data blocks for reuse. Although there are many types and manufacturers of flash memory, the most common type is known as NOR flash. NOR flash, such as that available from Intel Corporation™, operates in the following fashion: Erased state is 1, programmed state is 0, a 0 cannot be changed back to a 1 except by an erase, and an erase must occur on a full erase block. For additional detail, the following document may be referenced: “Understanding the Flash Translation Layer (FTL) Specification”, Intel, 1998, and is incorporated herein fully by reference.\nThe emergence of commodity PCIe flash marks a remarkable shift in storage hardware, introducing a three-order-of-magnitude performance improvement over traditional mechanical disks in a single release cycle. PCIe flash provides a thousand times more random IOPS than mechanical disks (and 100 times more than SAS/SATA SSDs) at a fraction of the per-IOP cost and power consumption. However, its high per-capacity cost makes it unsuitable as a drop-in replacement for mechanical disks in all cases. As such, systems that have either or both the high performance of flash with the cheap capacity of magnetic disks in order to optimize these balancing concerns may become desired. In such systems, the question of how to arrange data across both such media, but also within such media is helpful in optimizing the requirements for wide sets of data. For example, in the time that a single request can be served from disk, thousands of requests can be served from flash. Worse still, IO dependencies on requests served from disk could potentially stall deep request pipelines, significantly impacting overall performance.\nThis background information is provided to reveal information believed by the applicant to be of possible relevance. No admission is necessarily intended, nor should be construed, that any of the preceding information constitutes prior art."}
{"text": "Cable hangers are commonly used to secure cables to structural members of antenna towers and or along tunnel walls. Generally, each cable is attached to a structural member by cable hangers mounted at periodically-spaced attachment points.\nAntenna towers and/or tunnels may be crowded due to the large numbers of cables required for signal-carrying. Over time, as systems are added, upgraded and/or expanded, installation of additional cables may be required. To conserve space, it may be desirable for each set of cable hangers to secure more than a single cable. Certain cable hangers have been constructed to secure multiple cables; other cable hangers have a stackable construction that permits multiple cable hangers to be interlocked extending outwardly from each mounting point/structural member. Stacked and multiple-cable-type cable hangers significantly increase the number of cables mountable to a single attachment point.\nOne popular stackable cable hanger is discussed in U.S. Pat. No. 8,191,836 to Korczak, the disclosure of which is hereby incorporated herein by reference in its entirety. One such cable hanger, designated broadly at 10, is shown in FIGS. 1 and 2. The hanger 10 includes curved arms 5 that extend from a flat base 6. Locking projections 7 extend from the free ends of the arms 5. As can be seen in FIGS. 1 and 2, the locking projections 7 are inserted into a reinforced hole 8 in a tower structure 4 to mount the hanger 10 thereon. The base 6 of the hanger 10 includes a reinforced hole 9 that can receive the projections of another hanger 10 to mount a second cable.\nAs can be best seen in FIG. 2, the arms 5 include arcuate sections 14 that together generally define a circle within which a cable can reside. Two cantilevered tabs 12 extend radially inwardly and toward the base 6 at one end of the arcuate sections 14, and two cantilevered tabs 16 extend radially inwardly and toward the base 6 from the opposite ends of the arcuate sections 14. The cantilevered tabs 12, 16 are deployed to deflect radially outwardly when the hanger 10 receives a cable for mounting; this deflection generates a radially inward force from each tab 12, 16 that grips the jacket of the cable.\nHangers can be “stacked” onto each other by inserting the locking projections 7 of one hanger into the large hole 9 of the next hanger. One variety of cable hanger of this type is the SNAP-STAK® hanger, available from CommScope, Inc. (Joliet, Ill.).\nThe SNAP-STAK® hanger is offered in multiple sizes that correspond to the outer diameters of different cables. This arrangement has been suitable for use with RF coaxial cables, which tend to be manufactured in only a few different outer diameters; however, the arrangement has been less desirable for fiber optic cables, which tend to be manufactured in a much greater variety of diameters. Moreover, fiber optic cables tend to be much heavier than coaxial cables (sometimes as much as three times heavier per unit foot), which induces greater load and stress on the hangers.\nMultiple approaches to addressing this issue are offered in co-assigned and co-pending U.S. Patent Publication No. 2016/0281881 to Vaccaro, the disclosure of which is hereby incorporated herein by reference in full. One cable hanger discussed in this publication is shown in FIGS. 3 and 4 and designated broadly at 610 therein. The cable hanger 610 is somewhat similar to the cable hanger 10, inasmuch as it has a base 606, curved arms 605 and locking projections 607 that resemble those of the hanger 10 discussed above. However, the cable hanger 610 also has flex members 618 that define chords across the arcuate sections 614 of the arms 605. As can be seen in FIG. 4, cantilevered gripping members 612, 616 extend from the flex members 618 and into the cable-gripping space S within the arms 605. It can also be seen in FIG. 3 that the flex members 618 are tripartite, with two vertically offset horizontal runs 618a, 618c merging with the arcuate sections 614 of the arms 605 and a vertical run 618b extending between the horizontal runs 618a, 618c. The gripping members 612, 616 extend from opposite sides of the vertical run 618b and are vertically offset from each other.\nIn use, the cable hanger 610 is employed in the same manner as the cable hanger 10; a cable is inserted into the space S between the arms 605, which are then closed around the cable as the locking projections 607 are inserted into a mounting hole. The cantilevered gripping members 612, 616 can help to grip and to center the cable within the space S. The presence of the flex members 618, which are fixed end beams rather than cantilevered tabs, can provide additional gripping force beyond that of the cable hanger 10.\nIn view of the foregoing, it may be desirable to provide additional configurations of cable hangers to enable a technician to adapt to different cable sizes and mounting conditions."}
{"text": "1. Field of the Invention\nThis invention relates to a catalyst component or catalyst that is useful for the stereoregular polymerization or copolymerization of alpha-olefins and more particularly concerns a magnesium-containing supported titanium-containing catalyst component or catalyst that is useful for producing a homopolymer or copolymer of an alpha-olefin.\n2. Discussion of the Prior Art\nAlthough many polymerization and copolymerization processes and catalyst systems have been described for polymerizing or copolymerizing alpha-olefins, it is highly desirable to develop a catalyst component or a catalyst that has improved activity for catalyzing such reactions. It is also advantageous to tailor a process and catalyst system to obtain a specific set of properties of a resulting polymer or copolymer product. For example, in certain applications a product with a broader molecular weight distribution is desirable. Such a product has a lower melt viscosity at high shear rates than a product with a narrower molecular weight distribution. Many polymer or copolymer fabrication processes which operate with high shear rates, such as injection molding, oriented film, and thermobonded fibers, would benefit with a lower viscosity product by improving throughput rates and reducing energy costs. Thus, it is highly desirable to develop a catalyst or catalyst component that is useful for producing a homopolymer or copolymer of an alpha-olefin having a broadened molecular weight distribution. Also important is maintaining high activity and low atactic levels such as measured by hexane soluble and extractable materials formed during polymerization or copolymerization.\nMagnesium-containing supported titanium halide-based alpha-olefin polymerization or copolymerization catalyst components or catalyst systems containing such components are now well known in the art. Typically, these catalyst components and catalyst systems are recognized for their performance based on activity and stereospecificity. Numerous individual processes or process steps have been disclosed which have as their purpose the provision of improved supported, magnesium-containing, titanium-containing, electron donor-containing olefin polymerization or copolymerization catalysts. More particularly, Arzoumanidis et al., U.S. Pat. Nos. 4,866,022; 4,988,656; and 5,013,702 disclose a method for forming a particularly advantageous alpha-olefin polymerization or copolymerization catalyst or catalyst component that involves a specific sequence of specific individual process steps such that the resulting catalyst or catalyst component has exceptionally high activity and stereospecificity combined with very good morphology. A solid hydrocarbon-insoluble, alpha-olefin polymerization or copolymerization catalyst or catalyst component with superior activity, stereospecificity and morphology characteristics is disclosed as comprising the product formed by 1) forming a solution of a magnesium-containing species from a magnesium hydrocarbyl carbonate or magnesium carboxylate; 2) precipitating solid particles from such magnesium-containing solution by treatment with a transition metal halide and an organosilane as a morphology controlling agent; 3) reprecipitating such solid particles from a mixture containing a cyclic ether; and 4) treating the reprecipitated particles with a transition metal compound and an electron donor.\nArzoumanidis et al., U.S. Pat. No. 4,540,679 discloses a process for the preparation of a magnesium hydrocarbyl carbonate by reacting a suspension of a magnesium alcoholate in an alcohol with carbon dioxide and reacting the magnesium hydrocarbyl carbonate with a transition metal component. Arzoumanidis et al., U.S. Pat. No. 4,612,299 discloses a process for the preparation of a magnesium carboxylate by reacting a solution of a hydrocarbyl magnesium compound with carbon dioxide to precipitate a magnesium carboxylate and reacting the magnesium carboxylate with a transition metal component.\nWhile each of the processes of the aforesaid U.S. Pat. Nos. 4,540,679; 4,612,299; 4,866,022; 4,988,656; and 5,013,702 affords alpha-olefin polymerization or copolymerization catalysts or catalyst components which have high activity for polymerizing or copolymerizing alpha-olefins to produce homopolymer or copolymer products which have desirable characteristics, it is highly desirable to develop additional alpha-olefin polymerization or copolymerization catalysts or catalyst components--and methods for the manufacture thereof--that have even further improved catalytic activity and that afford polymers or copolymers which also have broadened molecular weight distribution.\nFor example, Karayannis, Cohen and Ledermann, pending U.S. patent application Ser. No. 07/862,960, filed Apr. 3, 1992, now U.S. Pat. No. 5,227,354, disclose a solid, hydrocarbon-insoluble catalyst or catalyst component and a method of production thereof, which are based on the catalyst or catalyst components and methods of production thereof, respectively, of the aforesaid U.S. Pat. Nos. 4,540,679; 4,612,299; 4,866,022; 4,988,656; and 5,013,702, wherein the resulting catalyst or catalyst component is a product formed by: A. forming a solution of a magnesium-containing species in a liquid, wherein the magnesium-containing species is formed by reacting a magnesium-containing compound with carbon dioxide or sulfur dioxide; B. precipitating solid particles from the solution of the magnesium-containing species by treatment with a titanium halide; and D. treating the precipitated particles with a titanium compound and an electron donor; wherein the treated precipitated particles from Step D comprise magnesium and vanadium components, and wherein vanadium is introduced into at least one of (i) the aforesaid magnesium-containing species in Step A by reacting the magnesium-containing compound or species with a vanadium-containing compound or complex, or (ii) the aforesaid solid particles precipitated in Step B by treatment of the magnesium-containing species with a titanium halide and a vanadium-containing compound or complex; or (iii) the aforesaid precipitated particles treated in Step D by treatment of the precipitated particles with a titanium compound, an electron donor and a vanadium-containing compound or complex that is free of a halide component. Use of the catalyst or catalyst component disclosed in the aforesaid Karayannis, Cohen and Ledermann pending patent application for the polymerization or copolymerization of an alpha-olefin affords polymers or copolymers which have a broadened molecular weight distribution, but such pending patent application does not disclose a substantial increase in catalytic activity for such polymerization or copolymerization.\nSimilarly Tachibana et al., U.S. Pat. No. 5,084,429 discloses a catalyst for use in polymerization of olefins which comprises a carrier mainly composed of a magnesium compound precipitated from a solution and a catalytic component supported on the carrier and selected from titanium halides, vanadyl halides and vanadium halides is described. The catalyst is obtained by a process which comprises: (A) mixing (a) at least one magnesium compound with (c) a saturated or unsaturated monohydric or polyhydric alcohol for reaction in dissolved state in the presence of (b) carbon dioxide in an inert hydrocarbon solvent to obtain component (A); (B) subjecting the component (A) to mixing and reaction with (d) a titanium and/or a vanadyl halide and/or a vanadium halide of the general formula, VX.sub.r (OR.sup.8).sub.4-n and also with (e) at least one boron compound, silicon compound and/or siloxane compound thereby obtaining solid product (I): (C) reacting the solid product (I) with (f) a cyclic ether with or with R.sup.12 OH thereby causing dissolution and re-precipitation to obtain solid product (II): and (D) subjecting the solid product (II) to further reaction with (g) component (B) consisting of a titanium halide and/or a vanadyl halide and/or a vanadium halide and/or a SiX.sub.S (OR.sup.9).sub.4-s, thereby obtaining solid product (III), followed either by further reaction with a mixture of the component (B) and (h) an electron donor or by reaction of (g) with the solid product (III) obtained by the reaction between the solid product (II) and (h) or (h) with (j) electron donor, thereby obtaining solid product (IV) for use as the catalytic component.\nCatalysts for the polymerization of olefins containing other relevant combinations of metal components have also been disclosed. For example, Albizzati et al., U.S. Pat. No. 5,082,817 discloses a catalyst for the polymerization of olefins, obtained by means of the reaction of: (a) a compound of a transition metal, typically titanium, containing at least one metal-halogen linkage, supported on a magnesium halide in the active form, with (b) a compound of titanium, zirconium or hafnium containing at least one metal-carbon linkage. Similarly, Howard et al., U.S. Pat. No. 4,228,263 discloses a catalyst for the polymerization of propylene, which is the reaction product of (a) a metal oxide such as aluminum oxide, titanium oxide, silica and magnesia or physical mixtures thereof, and (b) an organometallic compound of zirconlure, titanium or hafnium.\nIn addition, polymer or copolymer morphology is often critical and typically depends upon catalyst morphology. Good polymer morphology generally involves uniformity of particle size and shape, a narrow particle size distribution, resistance to attrition and an acceptably high bulk density. Minimization of very small particles (fines) typically is very important especially in gas-phase polymerizations or copolymerizations in order to avoid transfer or recycle line pluggage. Therefore, it is highly desirable to develop alpha-olefin polymerization and copolymerization catalysts and catalyst components that have good morphology, and in particular, a narrow particle size distribution. Another property which is important commercially is the maintenance of an acceptably high bulk density."}
{"text": "Over the last decade there has been a critical need to provide remote users a native transmission control protocol/internet protocol (TCP/IP) link across ultra-high frequency satellite communications (UHF SATCOM), especially among military service branches. For many users (troops, small ships, submarines, aircraft and others) UHF SATCOM is the only available radio frequency (RF) path. While there have been attempts to provide products with such capability, these products are not interoperable with other products. Further, the products have limited capabilities and other drawbacks.\nTherefore, there is a current need for a protocol that may provide multi-user TCP/IP access across UHF satellite communications and that may operate on low bandwidth, handle long-latency communication, operate with half-duplex communications, and provide a high channel quality."}
{"text": "There exist in the marketplace today a number of different hook-fastener media to be described below. It is our belief that each of these existing hook-fasteners suffers from one or more shortcomings which hamper their utility and utilization."}
{"text": "The present invention relates to a colour ink jet printing method, to a set of inks, to ink jet printer cartridges, to substrates printed using the method and to an ink jet printer.\nInk jet printing is a non-impact printing technique which involves ejecting, thermally or by action of an oscillating piezo crystal, droplets of ink from one or more fine nozzles directly onto a substrate. The ink may be aqueous, solvent or hot melt based.\nThe printing of textiles is conventionally carried out by screen or roller printing using gravure engraved cylinders. The design to be printed has to be engraved on to a cylinder and each individual colour in a design requires the application of a separate screen with a colour premixed to the required shade. This is a long and slow process and it can take many months for a textile design to appear as a printed textile. Consequently, there is a demand for a printing process which enables new designs to be printed onto a textile quickly for proofing purposes and for small production runs.\nInk jet printing of textiles offers the potential to transfer a design to a textile much faster than traditional textile printing methods. However, conventional colour ink jet printers operate with a colouring system which uses three subtractive primary colours (Cyan, Magenta and Yellow) together with black, hereinafter referred to as CMYK. This colouring system only provides a limited range of colours compared with conventional textile printing methods.\nThe limited colour range provided by conventional colour ink jet printers also limits their applicability in other imaging technologies, particularly where high resolution images are required from digital cameras or when printing images from the Internet. These xe2x80x9cphotorealisticxe2x80x9d applications require a wide range of colours to produce images of the same photographic quality as those prepared by conventional screen printing methods.\nWe have surprisingly found that the colour gamut of an ink jet printer is extended close to that obtained using conventional printing methods when a combination of inks containing dyes with specific calorimetric positions in colour space are used together with a CMYK colouring system in an ink jet printer.\nAccording to a first aspect of the present invention there is provided a method for the coloration of a substrate comprising ink jet printing a first and second set of inks onto the substrate wherein:\n(a) the first set of inks consists of one or more inks each of which independently contains a colorant selected from yellow, magenta, cyan and optionally black; and\n(b) the second set of inks comprises one or more inks each of which independently contains a dye selected from:\na yellow dye of Formula (1) or salt thereof: \nwherein:\nX is a labile group or atom;\nR1 is alkyl or xe2x80x94NH2;\nA is xe2x80x94NR2R3, xe2x80x94OR2 or halogen;\nR2 is H or optionally substituted alkyl; and\nR3 is optionally substituted phenyl;\nan orange dye of Formula (2) or salt thereof: \nwherein:\nR4 is H or optionally substituted alkyl;\nX is as hereinbefore defined;\nB is xe2x80x94NR5R6, xe2x80x94OR5 or halogen; and\nR5 and R6 independently are H or optionally substituted alkyl;\na red dye of Formula (3) or salt thereof: \nwherein:\nX is as hereinbefore defined;\nE is xe2x80x94NR7R8, xe2x80x94OR7 or halogen;\neach R7 is independently H or alkyl; and\nR8 is optionally substituted aryl;\nand a blue dye of Formula (4) or salt thereof: \nwherein:\neach X independently is as hereinbefore defined;\neach L independently is an optionally substituted alkylene group;\neach W independently is halogen or a group of the formula xe2x80x94OR9 or xe2x80x94NR10R11;\nR9 is H or alkyl, R10 is H or optionally substituted alkyl; and\nR11 is an optionally substituted aryl group.\nFIG. 1 compares the colour gamut obtained using a first set of inks alone (i.e. cyan, magenta and yellow) with that obtained using a first and second set of inks according to the present invention. In FIG. 1 the dots indicate the colour gamut from the first and second sets of inks. The crosses show the colour gamut from the first set of inks alone.\nThe coloration method of the present invention achieves a wide colour gamut by applying, in any combination, the inks from the hereinbefore defined first and second sets of inks to a substrate using an ink jet printer.\nPreferably, the first set of inks consists of yellow, magenta and cyan inks and preferably also a black ink. Accordingly, in a preferred embodiment, the first set of inks consists of an ink containing a yellow colorant, an ink containing a magenta colorant, an ink containing a cyan colorant and an ink containing a black colorant.\nThe second set of inks comprises from one to four, preferably two to four and especially four inks each of which independently contains a dye selected from the hereinbefore defined dyes of the Formulae (1) to (4).\nIn a preferred embodiment, the second set of inks comprises an ink containing a dye of Formula (1), an ink containing a dye of Formula (2) an ink containing a dye of Formula (3) and an ink containing a dye of Formula (4).\nIn view of the above preferences, in an especially preferred embodiment of the present invention the first set of inks consists of an ink containing a yellow colorant, an ink containing a magenta colorant, an ink containing a cyan colorant and an ink containing a black colorant; and the second set of inks comprises an ink containing a dye of Formula (1), an ink containing a dye of Formula (2) an ink containing a dye of Formula (3) and an ink containing a dye of Formula (4).\nAs will be understood, the ink of the first set of inks which contains the yellow colorant is a different shade of yellow to that of the ink of the second set of inks which contains the yellow dye of Formula (1) so as to maximise the colour gamut available to the ink jet printer.\nIt is preferred that each of the inks in the first and second set of inks are a different colour.\nThe inks may be applied to the substrate in any combination, provided that at least one ink from each of the two sets is applied. The particular combination of inks selected from the first and second sets of inks will be chosen to achieve a desired colour at a specific position on the substrate. Therefore, a single ink from the first or second set of inks or any combination of inks from the first and second sets may be applied to a particular point on a substrate to achieve a specific colour at that point. By way of illustration, if a red corresponding to the colour of the red ink of the second set of inks was required at a particular point on the substrate, then that ink alone would be applied to that part of the substrate by the ink jet printer.\nIn the dyes of Formula (1), it is preferred that:\nR1 is C1-4-alkyl or xe2x80x94NH2, more preferably methyl, ethyl or xe2x80x94NH2 and especially xe2x80x94NH2;\nR2 is H or optionally substituted C1-6-alkyl, more preferably H or C1-4-alkyl and especially H, methyl and ethyl;\nR3 is phenyl substituted by xe2x80x94COOH ,xe2x80x94SO3H or xe2x80x94OH, more preferably phenyl substituted by xe2x80x94SO3H; and\nthe group A is a group of the formula xe2x80x94NR2R3, xe2x80x94OR2 or Cl, more preferably xe2x80x94NR2R3.\nThe dyes of Formula (1) may be prepared by methods analogous to those described in the art for other similar azo dyes, for example as described in EP 559 331A1, page 5, wherein the compound of formula (4) has the formula H-A.\nIn the dyes of Formula (2), it is preferred that:\nR4 is H or optionally substituted C1-4-alkyl, more preferably H, methyl or ethyl and especially methyl;\nR5 and R6 are preferably each independently H or C1-4-alkyl, more preferably H; and\nB is preferably xe2x80x94NR5R6, xe2x80x94OR6 or Cl, more preferably xe2x80x94NR5R6 and especially xe2x80x94NHR6.\nThe dyes of Formula (2) may be prepared by methods analogous to those described in the art for similar azo dyes, for example as described in GB 859,990, Example 1.\nPreferred dyes of Formula (3) are those in which:\nR7 is H or C1-4-alkyl, more preferably H, methyl or ethyl and especially H;\nR8 is optionally substituted phenyl or naphthyl and especially phenyl optionally substituted by C1-4-alkyl; and\nE is xe2x80x94NR7R8, xe2x80x94OR7 or Cl, more preferably xe2x80x94NR7R8 and especially xe2x80x94NHR8.\nThe dyes of Formula (3) may be prepared by methods analogous to those described in the art for other similar azo dyes, for example as described in GB 834,304, Example 3.\nPreferred dyes of the Formula (4) are those in which:\nR9 is H or C1-4-alkyl, especially H;\nR10 is H or optionally substituted C1-6-alkyl, more preferably H or C1-4-alkyl and especially H or methyl;\nR11 is an optionally substituted phenyl, more preferably phenyl substituted with carboxy, hydroxy or sulpho and especially phenyl substituted with one or preferably two sulpho groups;\neach L independently is an optionally substituted C1-6-alkylene group, more preferably a C1-6-alkylene group and especially an alkylene group of the formula xe2x80x94CmH2mxe2x80x94,\nwherein m is an integer from 1 to 6, preferably 2 or 3; and each W independently is xe2x80x94OR9, xe2x80x94NR10R11 or Cl, more preferably xe2x80x94NR10OR11 and especially xe2x80x94NHR11.\nThe dyes of Formula (4) may be prepared by methods analogous to those described in the prior art for similar triphenodioxazine dyes, for example as described in EP 576 123 Al on pages 6 to 8 and Example 1.\nThe cyan, yellow, magenta and black colorants present in the inks of the first set of inks may be selected from any of the colorants suitable for use in ink jet printers utilising the conventional CMYK colouring system for example, cyan, yellow, magenta and black pigments and dyes. Preferably the colorants present in the first set of inks are dyes, more preferably dyes which contain fibre reactive groups.\nWhen any of the first set of inks contain a fibre reactive dye, any of the cyan, magenta, yellow and black reactive dyes listed in the Colour Index are suitable. Preferred reactive dyes comprise of one or more chromophores and one or more fibre reactive groups. Preferred fibre reactive groups are cellulose reactive groups, more preferably, vinylsulphone, sulphatoethylsulphone, halopyrimidine and holotriazine groups. Preferred halotriazine groups are fluoro and chlorotriazine groups, more preferably mono-chlorotriazine groups. Examples of suitable chromophores include azo, phthalocyanine, triphenodioxazine, anthraquinone and formazan chromophores. Especially preferred reactive dyes suitable for use in the first set of inks comprise one or more azo, phthalocyanine, triphenodioxazine, anthraquinone and formazan chromophore and one or more mono-chlorotriazine group. It is preferred that the reactive dye is water soluble.\nSpecific examples of reactive dyes suitable for use as colorants in the first set of inks include for example. Colour Index (C.I.) Reactive Red 31, C.I. Reactive Blue 71 and C.I. Reactive Black 8. These dyes are available from Zeneca Ltd.\nIn a preferred embodiment, the cyan colorant in the first set of inks is a dye of the Formula (5), or salt thereof \nwherein:\nPc is a phthalocyanine nucleus;\nZ is xe2x80x94NR12R13, xe2x80x94OR13 or halogen; wherein\nR12 is H or optionally substituted alkyl;\nR13 is H, optionally substituted alkyl or optionally substituted aryl;\nX is as hereinbefore defined;\na is 1 to 3;\nb is 1 to 3; and\n(a+b)=4.\nIt is preferred that R12 is H or optionally substituted C1-4-alkyl, more preferably H or C1-4-alkyl and especially H.\nPreferably R13 is H, optionally substituted C1-6-alkyl or optionally substituted phenyl, more preferably H or optionally substituted C1-4-alkyl and especially methyl or ethyl.\nWhen Z is halogen it is preferably chloro.\nPreferably Z is a group of the formula xe2x80x94NR12R13 or xe2x80x94OR13, more preferably xe2x80x94NHR13 or xe2x80x94OR13 and especially methoxy and ethoxy.\nThe dyes of Formula (5) may be prepared by methods analogous to those described in the art for similar phthalocyanine compounds. For example a suitable method is described in GB 805,562, Example 12.\nIn a preferred embodiment the yellow colorant present in the first set of inks is a dye of Formula (6) or salt thereof: \nwherein:\nY is a halogen or a group of the formula xe2x80x94NR15R16 or xe2x80x94OR16, wherein R15 is H or optionally substituted C1-4-alkyl, and R16 is H optionally substituted alkyl or aryl;\nR14 is H, optionally substituted alkyl or aryl; and\neach X independently is as hereinbefore defined.\nPreferably R14 is H or optionally substituted C1-4-alkyl, more preferably C1-4-alkyl and especially methyl or ethyl.\nPreferably R15 is H or C1-4-alkyl, more preferably H, methyl or ethyl and especially H.\nPreferably R16 is H, optionally substituted C1-4-alkyl or optionally substituted phenyl, more preferably H, C1-4-alkyl or phenyl and especially H, methyl and ethyl.\nWhen Y is halogen it is preferably chloro.\nIt is preferred that Y is a group of the formula xe2x80x94NR15R16, more preferably xe2x80x94NHR16 and especially xe2x80x94NH2.\nThe dyes of Formula (6) may be prepared by methods analogous to those described in the art for similar azo dyes. For example a suitable method is described in GB 1,271,226, Examples 2 to 153.\nIn a preferred embodiment, the magenta colorant present in the first set of inks is a dye of the formula (7) or salt thereof: \nwherein:\nG is halogen or a group of the formula xe2x80x94NR17R18 or xe2x80x94OR17, wherein R17 is H or optionally substituted alkyl;\nR18 is H, optionally substituted aryl or optionally substituted alkyl;\nR19 and R20 are each independently C1-4-alkyl; and\nX is as hereinbefore defined.\nIt is preferred that R17 is H or optionally substituted C1-4-alkyl, more preferably H or C1-4-alkyl and especially H or methyl.\nPreferably R18 is H, optionally substituted phenyl or optionally substituted C1-6-alkyl, more preferably H, C1-6-alkyl or phenyl substituted by C1-4-alkyl, xe2x80x94COOH or xe2x80x94SO3H, especially C1-6-alkyl and phenyl substituted by methyl, ethyl or xe2x80x94SO3H and more especially phenyl substituted with xe2x80x94SO3H and methyl.\nR19 and R20 are preferably independently methyl, ethyl or iso-propyl.\nWhen G is halogen it is preferably chloro.\nPreferably G is a group of the formula xe2x80x94NR17R18, more preferably xe2x80x94NHR18.\nThe dyes of Formula (7) may be prepared using methods analogous to those described in the art for other azo dyes. For example a suitable method is described in W094/22961.\nA preferred black colorant present in the first set of inks is a 1:2 Chromium or Cobalt complex, or mixture thereof, of Formula (8) or salt thereof: \nwherein:\nT is halogen or a group of the formula xe2x80x94NR21R22 or xe2x80x94OR21, wherein\nR21 is H or optionally substituted alkyl or optionally substituted aryl;\nR22 is H or optionally substituted alkyl;\nM is Cr or Co; and\nX is as hereinbefore defined.\nPreferably R21 is H, optionally substituted C1-6-alkyl or optionally substituted phenyl, more preferably H or C1-4-alkyl and especially H, methyl or ethyl.\nPreferably R22 is H or C1-4-alkyl, more preferably H.\nWhen T is halogen it is preferably chloro.\nIt is preferred that T is a group of the formula xe2x80x94NR21R22, more preferably xe2x80x94NHR21.\nPreferably a mixture of Cobalt and Chromium complexes of the dyes of Formula (8) are present in the black ink. Preferred mixtures comprise the components:\n(a) from 50 to 95, more preferably from 55 to 80 and especially from 60 to 75 parts of the 1:2 chromium complex; and\n(b) from 5 to 50, more preferably 20 to 45 and especially 40 to 25 parts of the 1:2 cobalt complex, wherein all parts are by weight and the parts (a)+(b)=100.\nThe dye of Formula (8) and the complexes thereof may be prepared using methods analogous to those described in the art for similar azo dyes. For example, a suitable method is disclosed in GB 985,481, Examples 1 and 2.\nA further preferred black colorant suitable for use in the first set of inks is a dye of the Formula (9) or a salt thereof. \nwherein:\nX is as hereinbefore defined; and\nR23 and R24 are each independently H or optionally substituted alkyl.\nPreferably R23 and R24 are each independently H or C1-6-alkyl optionally substituted by xe2x80x94OH or SO3H, more preferably H or C1-4-alkyl and especially H.\nThe sulpho group on the phenyl ring is preferably attached at the para position relative the second sulpho group on the ring.\nIt is preferred that the dye of Formula (9) is mixed with a small quantity of a yellow and/or a red dye to give a neutral black shade.\nA preferred mixture of dye comprises:\n(a) from 50 to 95, more preferably from 60 to 80 parts of the dye of Formula (9);\n(b) from 5 to 20, more preferably from 5 to 15 parts of a yellow dye; and\n(c) from 10 to 30, more preferably from 15 to 25 parts of a red dye;\nwherein all parts are by weight and the sum of the parts (a), (b) and (c)=100. The red and yellow dyes which may be present in the mixture are preferably azo dyes, more preferably water-soluble azo dyes.\nA preferred red dye is the dye of the hereinbefore defined Formula (3) wherein X is Cl and E is 2-methylphenylamino. A preferred yellow dye is the dye of the hereinbefore defined Formula (1) wherein X is Cl, A is 3-sulpho-N-methylaniline or 4-sulpho-N-methyl anilino and R1 is amino.\nAn especially preferred yellow dye is a 1:1 mixture of the dyes comprising:\n(i) a dye of Formula (1) in which X is Cl, A is 3-sulpho-N-methylanilino and R1 is amino; and\n(ii) a dye of Formula (1) in which X is Cl, A is 4-sulpho-N-methylanilino and R1 is amino.\nIn view of the foregoing preferences an especially preferred black colorant comprises:\n(a) from 60 to 80 parts of the dye of Formula (10);\n(b) from 15 to 25 parts of a 1:1 mixture of dyes of Formula (11); and\n(c) from 5 to 15 parts of a red dye of Formula (12);\nwherein all parts are by weight and the sum of the parts (a), (b) and (c)=100: \nThe dye of Formula (9) may be prepared using conventional methods for the preparation of azo dyes. For example a suitable method comprises:\n(1) diazotising the aniline disulphonic acid of the formula: \n(2) diazotising 2-amino-4-nitrobenzenesulphonic acid;\n(3) coupling the product of stage (2) with 8-amino-1-naphthol-3,6-disulphonic acid under mildly alkaline conditions;\n(4) coupling the product of stage (3) with the product of stage (1) under acidic conditions;\n(5) reducing the nitro group on the product of stage (4), for example by hydrogenation over a palladium catalyst;\n(6) condensing the product of stage (5) with the triazine compound of the formula: \nand\n(7) condensing the product of stage (5) with a compound of the formula NHR23R24 \nwherein X, R23 and R24 are as hereinbefore defined.\nA further preferred magenta colorant suitable for use in the first set of inks is a dye of the Formula (13) or salt thereof: \nwherein:\nQ is halogen or a group of the formula xe2x80x94NR25R26 or xe2x80x94OR25, wherein\nR25 is H or optionally substituted alkyl;\nR26 is H, optionally substituted aryl or optionally substituted alkyl; and\nX is as hereinbefore defined.\nIt is preferred that R25 is H or optionally substituted C1-4-alkyl, more preferably H or C1-4-alkyl and especially H or methyl.\nPreferably R26 is H, optionally substituted phenyl or optionally substituted C1-6-alkyl, more preferably H, C1-6-alkyl or phenyl substituted by C1-4-alkyl, xe2x80x94COOH or xe2x80x94SO3H, especially C1-6-alkyl and phenyl substituted by methyl, ethyl or xe2x80x94SO3H and more especially phenyl substituted with xe2x80x94SO3H and xe2x80x94COOH.\nWhen Q is halogen it is preferably chloro.\nPreferably Q is a group of the formula xe2x80x94NR25R26, more preferably xe2x80x94NHR26.\nThe dyes of Formula (13) may be prepared using methods analogous to those described in the art for other azo dyes. For example a suitable method is described in GB 899,376.\nAlthough Formulae (1) to (13) show the dyes in their free acid form, it is intended that salts of the dyes are included within the scope of the present invention. Thus the inks may contain dyes in their free acid and/or salt forms.\nIn view of the foregoing preferences for the colorants for the first set of inks, in a particularly preferred embodiment, the first set of inks comprises an ink containing a dye of Formula (5), an ink containing a dye of Formula (6) an ink containing a dye of Formula (7) and an ink containing a dye of Formula (8) or a dye of Formula (9); and the second set of inks comprises an ink containing a dye of Formula (1), an ink containing a dye of Formula (2) an ink containing a dye of Formula (3) and an ink containing a dye of Formula (4).\nThe labile group or atom represented by X in the dyes of Formulae (1) to (9) is a group or atom which is bound directly to the triazine nucleus and which is readily displaced therefrom in mildly alkaline conditions. Preferred labile groups include for example, a sulphonic acid group; a thiocyano group; a quaternary ammonium group, for example a trialkyl ammonium group or an optionally substituted pyridinium group, for example 3- or 4-caboxypyridinium; xe2x80x94OSO3H; or CH3COOxe2x80x94.\nPreferred labile atoms include halogens, more preferably F, Cl or Br and especially Cl.\nWhen the substituents on the dyes of Formulae (1) to (13) are optionally substituted, preferred substituents are selected from xe2x80x94OH, xe2x80x94SO3H, xe2x80x94COOH, xe2x80x94CN, xe2x80x94NO2, xe2x80x94PO3H2, halogen, especially Cl or Br, C1-4-alkyl and C1-4-alkoxy.\nIt is preferred that the colorants/dyes present in the inks are purified prior to incorporation into the ink by removing substantially all of the inorganic salts and other by products which may be present in the colorant/dye. Suitable purification processes include reverse osmosis and/or ultrafiltration.\nThe inks used in the first and second set of inks of the present invention comprise the relevant colorant or dye and a medium. The inks may comprise a mixture of two or more different colorants or dyes, or salts thereof, to provide an ink composition of the desired colour for use as an ink of the first or second set of inks of the present invention.\nThe medium for the ink compositions may be a liquid or a low melting point solid. Liquid media may be aqueous or solvent-based.\nIt is preferred that when dyes are used in the inks they are dissolved completely in the aqueous or solvent medium to form a solution.\nPreferred inks are those comprising a dye, preferably a dye of the Formula (1) to (9) as hereinbefore defined and an aqueous medium.\nThe inks used in the method of the present invention preferably contain from 0.5% to 20%, more preferably from 0.5% to 15%, and especially from 2% to 12% by weight of the dye/colorant based on the total weight of the ink It is preferred that where the colorant is a dye, it has a solubility of around 10% or more to allow the preparation of concentrates which may be used to prepare more dilute inks and to minimise the chance of precipitation of the dye if evaporation of the liquid medium occurs during use of the ink.\nWhere the colorant is a dye its solubility can be enhanced by converting the sodium salt, in which form it is usually synthesised, either partially or wholly, into the lithium or ammonium salt Purification of the dye can be conveniently accomplished by use of membrane separation processes to separate unwanted by-products and inorganic materials from the solution or dispersion of the dye, followed by partial or complete exchange of the counter ion.\nWhere the liquid medium is aqueous based it is preferably water or a mixture of water and one or more water-soluble organic solvent. The weight ratio of water to organic solvent(s) is preferably from 99:1 to 1:99, more preferably from 95:1 to 50:50 and especially from 90:10 to 60:40.\nThe water-soluble organic solvent is preferably selected from C1-4-alkanols such as methanol, ethanol, n-propanol, isopropanol, n-butanol, sec-butanol, tert-butanol or isobutanol; amides such as dimethylformamide or dimethylacetamide; ketones or ketone-alcohols such as acetone or diacetone alcohol; ethers such as tetrahydrofuran or dioxane; oligo- or poly-alkylene glycols such as diethylene glycols, triethylene glycol, polyethylene glycol or polypropylene glycol; alkenyleneglycols or thioglycols containing a C2-C6-alkylene group such as ethylene glycol, propylene glycol, butylene glycol, pentylene glycol or hexylene glycol, thioglycol and thiodiglycol; polyols such as glycerol or 1,2,6-hexanetriol; C1-4-alkyl-ethers or polyhdric alcohols such as 2-methoxyethanol, 2-2-(2-methoxyethoxy)ethanol, 2-(2-ethoxyethoxy)-ethanol, 2-[2-(2-methoxyethoxy)ethoxy]ethanol, 2-[2-(2-ethoxyethoxy)-ethoxy]-ethanol; heterocyclic ketones, such as 2-pyrrolidone and N-methyl-2-pyrrolidone; or mixtures containing two or more of the aforementioned water-soluble organic solvents for example thiodiglycol and a second glycol or diethylene glycol and 2-pyrrolidone.\nPreferred water-soluble organic solvents are 2-pyrrolidone; N-methylpyrrolidone; alkylene glycols and oligo-alkylene glycols, such as ethylene glycol, diethylene glycol, triethylene glycol; and lower alkyl ethers of polyhydric alcohols such as 2-methoxy-2-ethoxy-2-ethoxy-ethanol; polyethylene glycols with a molecular weight of up to 500; and thioglycols such as thiodiglycols. A preferred specific solvent mixture is a binary or ternary mixture of water and diethylene glycol and/or, 2-pyrrolidone or N-methylpyrrolidone in weight ratios 75-95:25-5 and 60-98:1-20:1-20 respectively. An especially preferred specific solvent mixture is a binary or tertiary mixture of water and thiodiglycol and/or 2-pyrrolidone or N-methylpyrrolidone in weight ratios 75-98:25-2 and 60-90:5-20:5-20 respectively.\nExamples of suitable aqueous ink media are given in U.S. Pat. Nos. 4,963,189, 4,703,113, 4,626,284, EP 4,251,50A and U.S. Pat. No. 5,207,824.\nWhen aqueous inks are used in the present invention, they preferably also contain a humectant to inhibit evaporation of water and a preservative to inhibit the growth of fungi, bacteria and/or algae in the solution. Examples of suitable humectants are, propan-1,2-diol, butan-1,2-diol, butan-2,3-diol and butan-1,3-diol. However, the presence of small amounts, up to about 10%, preferably not more than 5%, in total, of polyols having two or more primary hydroxy and/or primary alcohols is acceptable, although the ink is preferably free from such compounds\nWhere the liquid medium is solvent based the solvent is preferably selected from ketones, alkanols, aliphatic hydrocarbons, esters, ethers, amides or mixtures thereof. Where an aliphatic hydrocarbon is used as the solvent a polar solvent such as an alcohol, ester, ether or amide is preferably added. Preferred solvents include ketones, especially methyl ethyl ketone and alkanols especially ethanol and n-propanol.\nSolvent based ink compositions are used where fast drying times are required and particularly when printing onto hydrophobic substrates such as plastics, metal or glass.\nWhere the medium for an ink composition is a low melting point solid the melting point of the solid is preferably in the range from 60xc2x0 C. to 125xc2x0 C. Suitable low melting point solids include long chain fatty acids or alcohols, preferably those with C18-24 chains, or sulphonamides. The dye or colorant may be dissolved in the low melting point solid or may be finely dispersed in it.\nThe inks may optionally contain other components conventionally used in inks for ink jet printing. For example, viscosity and surface tension modifiers, corrosion inhibitors, kogation reducing additives, surfactants and anti-cockle agents, for example those disclosed in U.S. Pat. No. 5,207,824, column 3, line 13 to column 4, line 21, which is included herein by reference thereto.\nIf desired, the inks may be buffered to a pH of from 5 to 8, more preferably from 6 to 7, with a suitable buffer such as the sodium salt of metanillic acid or an alkali metal phosphate.\nWhere the ink jet printing technique involves the charging and electrically-controlled deflection of drops, for example in a continuous ink jet printer, the inks preferably also contain a conducting material such as an ionised salt to enhance and stabilise the charge applied to the drops. Suitable salts for this purpose are alkali metal salts of mineral acids.\nIn view of the foregoing preferences, a particularly preferred ink composition is an aqueous ink comprising the components:\n(a) 0.5 to 20 parts of the relevant dye;\n(b) 50 to 98 parts water; and\n(c) 2 to 50 parts water soluble organic solvent;\nwherein all parts are parts by weight and the parts (a)+(b)+(c)=100. In addition to the components (a), (b) and (c), the ink may also contain further components as hereinbefore mentioned.\nThe ink jet printer used for the application of the inks to the substrate forms each ink into small droplets by ejection from a reservoir through a small orifice (the ink jet nozzle) so that the ink droplets are directed at the substrate during relative movement between the substrate and the reservoir. This process is commonly referred to as ink jet printing. The ink may be applied to the substrate using a xe2x80x9ccontinuousxe2x80x9d or a xe2x80x9cdrop on demandxe2x80x9d printer, both of which are well known in the art. Continuous ink jet printers produce a stream of ink from the ink jet nozzle which is formed into droplets and directed to the substrate via a suitable control means. Drop on demand printers eject individual droplets of ink from the ink jet nozzle in response to a control signal. Preferred drop on demand ink jet printers for use in the present invention are piezoelectric and thermal ink jet printers. In thermal ink jet printing, programmed pulses of heat are applied to the ink in the reservoir by means of a resistor adjacent to the ink jet nozzle. In piezo-electric printers ink droplets are ejected from the ink jet nozzle using a piezoelectric transducer. The transducer oscillates in response to an electrical control signal, thereby creating a pressure wave in the reservoir adjacent to the ink jet nozzle which ejects droplets of ink from the nozzle.\nInk jet printers suitable for use in the present invention apply at least one ink from the first set of inks and at least one ink from the second set of inks, wherein the first and second sets of inks are as hereinbefore defined.\nPreferably the ink jet printer applies a first set of inks consisting of black, cyan, magenta and yellow inks and from one to four, preferably from two to four and especially four of the hereinbefore defined inks of the second set of inks.\nInk jet printers suitable for use in the present invention require a means to access each of the inks comprising the first and second sets of inks, for example from separate ink tanks containing each ink.\nPreferably the inks comprising the first and second set of inks are contained in an ink cartridge suitable for use in an ink jet printer. The ink cartridge preferably comprises a container in which are held the first and second set of inks. Preferably the ink jet printer contains the ink cartridge, for example within a suitable housing in the printer.\nIt is preferred that the printer has a separate channel to direct each ink colour in the first and second set of inks to a nozzle, or array of nozzles, dedicated to that ink colour. The separate ink channels for each ink colour avoids cross contamination of inks in the ink jet head which could result in an undesirable colour on the printed substrate. Accordingly, preferred ink jet printers have from 5 to 8, more preferably from 6 to 8 and especially 8 ink channels. An example of an ink jet printer suitable for use in the present invention is described in EP 616 893 A2.\nThe substrate used in the ink-jet method of the present invention may be paper, plastics, textile, metal, ceramic or glass and is preferably paper, plastic or a textile material, especially a natural, semi-synthetic or synthetic material.\nExamples of natural textile materials include wool, silk, hair and cellulosic materials, particularly cotton, jute, hemp, flax and linen.\nExamples of synthetic and semi-synthetic materials include polyamides, polyesters, polyacrylonitriles and polyurethanes.\nPreferred substrates include overhead projector slides paper and textile materials. Preferred papers include plain and treated papers. Preferred textile materials are cellulosic materials such as cotton. Especially preferred substrates are treated papers suitable for high resolution xe2x80x9cphoto-realisticxe2x80x9d image printing and cellulosic textile materials.\nAccording to a second aspect of the present invention there is provided a set of ink jet printing inks comprising a first and second set of inks wherein:\n(a) the first set of inks consists of a yellow ink, a magenta ink, a cyan ink and optionally a black ink; and\n(b) the second set of inks comprises one or more inks selected from:\nan ink containing a dye of Formula (1) or salt thereof;\nan ink containing a dye of the Formula (2) or a salt thereof;\nan ink containing a dye of Formula (3) or a salt thereof; and\nan ink containing a dye of Formula (4) or a salt thereof;\nwherein the dyes of Formulae (1), (2), (3) and (4) are as hereinbefore defined in the first aspect of the present invention.\nPreferred inks in the set according to the second aspect of the present invention are as defined in the first aspect of the present invention.\nAccordingly a preferred set of ink jet printing inks comprises:\n(a) a first set of inks consisting of:\na cyan ink containing a dye of the Formula (5);\na yellow ink containing a dye of Formula (6);\na magenta ink containing a dye of the Formula (7) or Formula (13); and\na black ink containing a dye of the Formula (8) or Formula (9); and\n(b) a second set of inks comprising from one to four, preferably from two to four and especially four inks selected from:\nan ink containing a dye of Formula (1);\nan ink containing a dye of Formula (2);\nan ink containing a dye of Formula (3);and\nan ink containing a dye of Formula (4).\nAccording to a third aspect of the present invention there is provided a paper or an overhead projector slide, a metal, glass or ceramic substrate or textile material coloured by means of the method according to the first aspect of the present invention.\nAccording to a fourth aspect of the present invention there is provided a process for the coloration of a textile material using ink-jet printing which comprises the steps:\ni) applying to the textile material by ink-jet printing the inks from the first and second sets of ink in accordance with the method of the first aspect of the present invention; and\nii) heating the textile material at a temperature from 50xc2x0 C. to 250xc2x0 C. to fix the dyes and colorants on the material.\nThe process for coloration of a textile material by ink-jet printing preferably comprises a pre-treatment of the textile material with an aqueous pretreatment composition comprising a water-soluble base, a hydrotropic agent and a thickening agent followed by removing water from the pre-treated textile material to give a dry pre-treated textile material which is subjected to ink-jet printing in step i) above.\nThe pre-treatment composition preferably comprises a solution of the base and the hydrotropic agent in water containing the thickening agent.\nThe base is preferably an inorganic alkaline base, especially a salt of an alkali metal with a weak acid such as an alkali metal carbonate, bicarbonate or silicate or an alkali metal hydroxide. The amount of base may be varied within wide limits provided sufficient base is retained on the textile material after pre-treatment to promote the dyeing of the pre-treated textile material. Where the base is sodium bicarbonate it is convenient to use a concentration of from 1% to 5% by weight based on the total weight of the composition.\nThe hydrotropic agent is present to provide sufficient water to promote the fixation reaction between the dye and the textile material during the heat treatment, in step (ii) above, and any suitable hydrotropic agent may be employed. Preferred hydrotropic agents are urea, thiourea and dicyandiamide. The amount of hydrotropic agent depends to some extent on the type of heat treatment. If steam is used for the heat treatment generally less hydrotropic agent is required than if the heat treatment is dry, because the steam provides a humid environment. The amount of hydrotropic agent required is generally from 2.5% to 50% by weight of the total composition with from 2.5% to 10% being more suitable for a steam heat treatment and from 20% to 40% being more suitable for a dry heat treatment.\nThe thickening agent may be any thickening agent suitable for use in the preparation of print pastes for the conventional printing of cellulose reactive dyes. Suitable thickening agents include alginates, especially sodium alginate, xantham gums, monogalactam thickeners and cellulosic thickeners. The amount of the thickening agent can vary within wide limits depending on the relationship between concentration and viscosity. However, sufficient agent is preferred to give a viscosity from 10 to 1000 mPa.s, preferably from 10 to 100 mPa.s, (measured on a Brookfield RVF Viscometer). For an alginate thickener this range can be provided by using from 10% to 20% by weight based on the total weight of the pretreatment composition.\nThe remainder of the pre-treatment composition is preferably water, but other ingredients may be added to aid fixation of the dye to the textile material or to enhance the clarity of print by inhibiting the diffusion (migration) of dye from coloured areas to non-coloured areas before fixation.\nExamples of fixation enhancing agents include cationic polymers and quaternary ammonium compounds. Suitable cationic polymers include for example, a 50% aqueous solution of a dicyanamide/phenol formaldehydelammonium chloride condensate e.g. MATEXIL FC-PN (available from ICI), which have a strong affinity for the textile material and the dye and thus increase the fixation of the dye on the textile material. Suitable quaternary ammonium compounds include for example those described in our EP 534 660 A1, incorporated herein by reference thereto, such as distearyl dimethylammonium chloride.\nExamples of anti-migration agents are low molecular weight acrylic resins, e.g. polyacrylates, such as poly(acrylic acid) and poly(vinyl acrylate).\nIn the pre-treatment stage of the present process the pre-treatment composition is preferably evenly applied to the textile material. Where a deeply penetrated print or a deep shade is required the pretreatment composition is preferably applied by a padding or similar process so that it is evenly distributed throughout the material. However, where only a superficial print is required the pre-treatment composition can be applied to the surface of the textile material by a printing procedure, such as screen or roller printing, ink jet printing or bar application.\nIn the pre-treatment stage of the present process, water may be removed from the pre-treated textile material by any suitable drying procedure such as by exposure to hot air or direct heating, e.g. by infra-red radiation, or micro-wave radiation, preferably so that the temperature of the material does not exceed 100xc2x0 C.\nThe application of the ink composition to the textile material, stage (i) of the present process, is as hereinbefore defined for the coloration method of the first aspect of the present invention. It is preferred that each ink of the first set of inks contains a reactive dye.\nAfter application of the ink, it is generally desirable to remove water from the printed textile material at relatively low temperatures ( less than 100xc2x0 C.) prior to the heat applied to fix the dye on the textile material as this has been found to minimise the diffusion of the dye from printed to non-printed regions. As with the pretreated textile material removal of water is preferably by heat, such as by exposure to hot air or to infra-red or micro-wave radiation.\nIn stage (ii) of the present process, the printed textile material is submitted to a short heat treatment, preferably after removal of water by low-temperature drying, at a temperature from 100xc2x0 C. to 200xc2x0 C. by exposure to dry or steam heat for a period of up to 20 minutes. If a steam (wet) heat treatment is used, the printed material is preferably maintained at 100-105xc2x0 C. for from 5 to 15 minutes whereas if a dry heat treatment is employed the printed material is preferably maintained at 140-160xc2x0 C. for from 2 to 8 minutes.\nAfter allowing the textile material to cool, unfixed dye and other ingredients of the pretreatment and dye compositions may be removed from the textile material by a washing sequence, involving a series of hot and cold washes in water and aqueous detergent solutions before the textile material is dried.\nAccording to a fifth aspect of the present invention there are provided textile materials, especially cellulosic textile materials, coloured by means of the process according to the fourth aspect of present invention.\nAccording to a sixth aspect of the present invention there is provided a photorealistic print prepared using the coloration method according to the first aspect of the present invention.\nA photorealistic print is a high resolution print which reproduces an image from digital data corresponding to that image. The photorealistic prints prepared in accordance with the present invention exhibit excellent colour yield and a wide colour gamut compared with images formed by conventional CMYK ink jet printers.\nExamples of systems which prepare images in digitised form include digital cameras, images stored on a CD ROM system, images that have been digitised using an optical scanning device and computer graphics programs such as CAD/CAM systems.\nThe digital data corresponding to the image is used by a computer to control the discharge of inks selected from the first and second sets of ink from an ink jet printer in accordance with the coloration method of the first aspect of the present invention, thereby providing a photorealistic print corresponding to the digitised image on a substrate.\nThe substrate used in the preparation of photorealistic prints may be any of the hereinbefore mentioned substrates. However, preferred substrates are paper, more preferably coated paper. Examples of paper and coated paper substrates suitable for use in the preparation of photorealistic prints include Hewlett Packard coated papers such as HP 516347, HP Premium Coated Paper and HP Photopaper, Stylus Pro 720 dpi Coated Paper, Epson Photo Quality Glossy Film (available from Seiko Epson Corp.), Epson Photo Quality Glossy Paper (available from Seiko Epson Corp.) Canon HR 101 High Resolution Paper (available from Canon), Canon GP 201 Glossy Paper (available from Canon), and Canon HG 101 High Gloss Film (available from Canon)."}
{"text": "1. Field of the Invention\nThe present invention relates to a photographic lens which can focus from infinity to a very close distance. In particular, the present invention is directed to such photographic lens which can be used at 1:1 magnification at maximum.\n2. Description of the Prior Art\nLenses hitherto known useful for photography at a very short object distance are classified into two groups.\n(1) Lenses of the type comprising a common photographic lens and a supplementary lens for close-up photography (so-called close-up lens) mounted on the common photographic lens; PA1 (2) Lenses of the type which are particularly designed for short distance photography. This type of lens is, when used, shifted forward as a unit from the film plane.\nThese known systems for photographing a very close object have the following disadvantages in any case:\nIn the case of the first mentioned system it is impossible to continuously change the magnification at which photography is carried out. For photography within a desired range of magnification, a number of supplementary lenses are required. In addition, exchange of such supplementary lenses is very troublesome to the user.\nIn the case of the latter mentioned system, it is possible to continuously change the magnification. However, a large and complicated mechanism is needed for shifting the lens. If portability is a consideration, the range of focusing movement available for the lens is limited. For the known system, the available photographic magnification is generally in the order of .times.0.5 at most.\nFurthermore, it is desired that a close-up photographic lens satisfy various requirements at the same time. The lens should not only be able to take a photography approaching the object but also to take a photograph at a sufficiently large magnification while keeping a sufficiently large distance from the object to the lens (working distance). The lens should also be as small as possible in size. All of the lenses according to the prior art can not fully satisfy these requirements at the same time."}
{"text": "The detection of pathogenic microorganisms in biological fluids should be performed in the shortest possible time, in particular in the case of septicemia for which the mortality remains high in spite of the broad range of antibiotics which are available to doctors. The presence of biologically active agents such as a microorganism in a patient's body fluid, especially blood, is generally determined using blood culture bottles. A small quantity of blood is injected through an enclosing rubber septum into a sterile bottle containing a culture medium, and the bottle is then incubated at 37° C. and monitored for microorganism growth.\nInstruments currently exist on the market in the U.S. that detect the growth of a microorganism in a biological sample. One such instrument is the BacT/ALERT® 3D instrument of the present assignee bioMérieux, Inc. The instrument receives a blood culture bottle containing a blood sample, e.g., from a human patient. The instrument incubates the bottle and periodically during incubation an optical detection unit in the incubator analyzes a colorimetric sensor incorporated into the bottle to detect whether microbial growth has occurred within the bottle. The optical detection unit, bottles and sensors are described in the patent literature, see U.S. Pat. Nos. 4,945,060; 5,094,955; 5,162,229; 5,164,796; 5,217,876; 5,795,773; and 5,856,175, the entire content of each of which is incorporated by reference herein. Other prior art of interest relating generally to the detection of microorganisms in a biological sample includes the following patents: U.S. Pat. Nos. 5,770,394, 5,518,923; 5,498,543, 5,432,061, 5,371,016, 5,397,709, 5,344,417 and its continuation U.S. Pat. Nos. 5,374,264, 6,709,857; and 7,211,430, the entire content of each of which is incorporated by reference herein.\nSubstantial, and potentially life saving, clinical benefits for a patient are possible if the time it takes for detection of a microbial agent in a blood sample and reporting the results to a clinician could be reduced. A system that meets this need has heretofore eluded the art. However, such rapid detection of a microbial agent in a biological sample such as a blood sample is made possible by apparatus described herein.\nThe disclosed system and methods combines a detection system operative to detect a container containing a test sample (e.g., a biological sample) as being positive for microbial agent presence. The systems and methods of this disclosure have the potential to: (a) reduce laboratory labor and user errors; (b) improve sample tracking, traceability and information management; (c) interface to laboratory automation systems; (d) improve work-flow and ergonomics; (e) deliver clinically relevant information; (f) faster results.\nMany further advantages and benefits over the prior art will be explained below in the following detailed description."}
{"text": "1. Field of the Invention\nThe present invention relates to a hydraulic cylinder unit used as an actuator, and more particularly, to a simplified structural hydraulic cylinder providing simple maintenance and commonly used parts.\n2. Description of the Related Art\nThe prior art shows a hydraulic cylinder unit 51 of a type, as shown in FIG. 4, in which oil is pumped from one end side of a cylinder to two oil chambers defined in the cylinder by a piston.\nThe hydraulic cylinder unit 51 comprises a cylinder body 52 including an outer cylinder 53, an inner cylinder 54, a lower cap 55 and an upper cap 56; and a piston 57 which is slidable in the inner cylinder 54. Oil is pumped from the side of the lower cap 55 to an upper chamber a and a lower chamber b (the piston 57 is at the maximum compression position, and b is indicated by a line) defined in the inner cylinder 54 by the piston 57.\nThat is, the lower cap 55 is provided with an upper chamber port 55a and a lower chamber port 55b. The upper chamber port 55a is communicated with an oil passage c formed between the outer and inner cylinders 53 and 54 through a first communication passage R.sub.1, and the oil passage c is in communication with the upper chamber a through a communication passage Q of the upper cap 56.\nThe lower chamber port 55b is in communication with the lower chamber b through a second communication passage R.sub.2.\nIn such a hydraulic cylinder unit 51, in order to move the piston 57 upward, a hydraulic oil is supplied from the lower chamber port 55b to increase the hydraulic pressure in the lower chamber b, and oil in the upper chamber a is discharged from the upper chamber port 55a, and in order to move the piston 57 downward on the contrary, hydraulic oil is supplied from the upper chamber port a, and oil in the lower chamber b is discharged from the lower chamber port 55b.\nHowever, in this conventional structure, since the number of constituent parts of the cylinder body 52 is high, it takes much time to assemble the device. Further, since the upper chamber port 55a, the lower chamber port 55b, the first communication passage R.sub.1 and the second communication passage R.sub.2 are all integrally formed in the lower cap 55, there are problems due to complications in processing of the passages, the number of processing steps is high, the maintenance operation requires excessive time and labor, the cost is increased, and the lower cap is prone to be large in size.\nFurther, since a joint is threaded into the port portion, the port portion may be damaged if it is repeatedly attached and detached many times. In such a case, it is necessary to exchange the entire lower cap.\nFurthermore, if it is necessary to vary a diameter of the inner cylinder, and to reduce a diameter of the piston, the entire inner cylinder must be integrally exchanged, and it is difficult to commonly use the parts."}
{"text": "Currently, hoppers and powered screw conveyors are used for delivering feed to animals. Such feed systems, however, have been known for being problematic and requiring periodic maintenance and attention in order to keep the feed flowing and the systems operating."}
{"text": "1. Field of the Invention\nThe present invention is directed generally to a cosmetic application system and method, and more particularly to a system and method for removing excess mascara from a mascara brush upon withdrawal from a container.\n2. Description of the Related Art\nVarious techniques and structures have been used to reduce the amount of mascara on a mascara brush upon removal from a container. However, a number of disadvantages associated with these techniques and structures has inhibited their widespread use and manufacture.\nIn particular, U.S. Pat. Nos. 4,194,848, 4,332,494, 4,407,311, 4,609,300 and 4,705,053 are directed to mascara applicators having a complex structure for varying the amount of mascara remaining on a brush after removal from a container. A flexible member is disposed in the neck of the container to provide some degree of variation in the amount of mascara removed from a brush as it passes through an opening in the container. However, each of these patents is directed to a complex structure, which is difficult and costly to manufacture. Moreover, many of these structures do not facilitate continuous variation of the amount of mascara to be removed from a brush. In addition, because these structures apply an equal force against the brush during removal and re-insertion of the brush into the container, these systems unnecessarily impede a user's ability to reinsert the brush into the container after each use. In U.S. Pat. No. 5,397,193, an additional attempt was made to provide a system for removing excess mascara from an applicator brush. In this system, a plurality of internal flexible bristles are used to remove excess mascara from the applicator brush. As with the aforementioned patents, this system is also costly and difficult to manufacture, and does not facilitate continuous variation in the amount of force to be applied to the mascara brush upon removal from its container. In addition, this system also unnecessarily impedes a user's ability to reinsert the brush into the container after each use.\nIn addition, each of the aforementioned systems, because of their complicated internal structure, is particularly difficult to clean. Accordingly, these go systems do not lend themselves for use with any form of reusable or interchangeable mascara system."}
{"text": "The present disclosure pertains to controllers and particularly to controllers having network switches."}
{"text": "The invention relates to scintillator materials, to a manufacturing process for obtaining them and to the use of said materials, especially in gamma-ray and/or X-ray detectors.\nScintillator materials are widely used in detectors for gamma rays, X-rays, cosmic rays and particles having an energy of the order of 1 keV and also above this value.\nA scintillator material is a material that is transparent in the scintillation wavelength range, which responds to incident radiation by emitting a light pulse.\nIt is possible to manufacture from such materials, which are generally single crystals, detectors in which the light emitted by the crystal that the detector contains is coupled to a light detection means and produces an electrical signal proportional to the number of light pulses received and to their intensity. Such detectors are used in particular in industry to measure thickness and grammage or coating weight, and in the fields of nuclear medicine, physics, chemistry and oil research.\nOne family of known scintillator crystals that is used is that of cerium-doped lutetium silicates. Cerium-doped Lu2SiO5 is disclosed in U.S. Pat. No. 4,958,080, and the U.S. Pat. No. 6,624,420 discloses Ce2x(Lu1-yYy)2(1-x)SiO5. Finally, U.S. Pat. No. 6,437,336 relates to compositions of the Lu2(1-x)M2xSi2O7 type, where M is at least partly cerium. These various scintillator compositions all have in common a high stopping power for high-energy radiation and give rise to intense light emission with very rapid light pulses.\nA desirable additional property is to reduce the amount of light emitted after the incident radiation stops (i.e. delayed luminescence or afterglow). Physically, this phenomenon, well known to those skilled in the art, is explained by the presence of electron traps in the crystallographic structure of the material. The phenomenon of scintillation relies on the photoelectric effect, which creates an electron-hole pair in the scintillator material. Upon recombination on an active site (a Ce3+ site in the aforementioned scintillators), the electron emits photons via a process that generally takes place in much less than one microsecond. The aforementioned scintillators, which are particularly rapid, result in a pulse duration that decreases with a first-order exponential constant of around 40 ns. However, the trapped electrons do not generate light, but their detrapping by thermal excitation (including at room temperature) gives rise to photon emission—the afterglow—, which still remains measurable after times of greater than one second.\nThis phenomenon may be unacceptable in applications in which it is desired to isolate each pulse, using very short windowing. This is particularly the case with CT (computed tomography) applications (scanners) that are well known in the medical or industrial sectors. When the CT system is coupled to a PET (Positron Emission Tomography) scanner, which is becoming the standard in industry, the poorer resolution of the CT affects the performance of the entire system and therefore the capability of the clinician to interpret the result of the complete PET/CT system. Afterglow is known to be completely unacceptable for these applications.\nCompositions of the lutetium silicates type, disclosed in U.S. Pat. No. 4,958,080 (of the LSO:Ce type, using the notation of those skilled in the art) and U.S. Pat. No. 6,624,420 (of the LYSO:Ce type) are known to generate a significant afterglow. In contrast, the compositions disclosed in U.S. Pat. No. 6,437,336 (of the LPS:Ce type) have the advantage of a much weaker afterglow. These results are given for example by L. Pidol, A. Kahn-Harari, B. Viana, B. Ferrand, P. Dorenbos, J. de Haas, C. W. E. van Eijk and E. Virey in “Scintillation properties of Lu2Si2O7:Ce3+, a fast and dense scintillator crystal”, Journal of Physics: Condensed Matter, 2003, 15, 2091-2102. The curve shown in FIG. 1 is extracted from this article and represents the amount of light detected in the form of the number of events (or counts) per mg of scintillator material as a function of time, under X-ray excitation for a few hours. The LPS:Ce composition gives a significantly better result in terms of afterglow.\nThe behavior of LYSO is very similar to that of LSO from this standpoint. The reduction in this afterglow forms the subject of the present application.\nThe afterglow property may be demonstrated more fundamentally by thermoluminescence (see S. W. S. McKeever “Thermoluminescence of solids”, Cambridge University Press (1985)). This characterization consists in thermally exciting a specimen after irradiation and measuring the light emission. A light peak close to room temperature at 300 K corresponds to an afterglow of greater or lesser magnitude depending on its intensity (detrapping). A peak at a higher temperature corresponds to the existence of traps that are deeper but less susceptible to thermal excitation at room temperature. This is illustrated in FIG. 2, extracted from the aforementioned article by L. Pidol et al., which shows, in another way, the advantage of a composition of the LPS type in terms of afterglow.\nHowever, compositions of the LPS type have the drawback of a lower stopping power than those of the LSO or LYSO type. This situation stems simply from the average atomic number of the compound and from the density of the associated phase."}
{"text": "1. Field of the Art\nThe present invention relates to an electronic typewriter having a carriage which is automatically returned upon entering of predetermined carriage-return data such as space data or hyphen data through the corresponding key on a keyboard, if the carriage-return data is entered to be executed within a predetermined automatic-carriage-return zone which consists of a desired number of columns just before or up to the right-hand margin of a line of printing.\n2. Related Art Statement\nIn making documents by means of an electronic typewriter, it is a common practice that a pair of successive space data are entered following period data at the end of a sentence in order to distinguish a space between sentences from a space between words of a sentence.\nIf the above typing method is used on an electronic typewriter having a function of executing an automatic carriage return, more specifically, if two sets of space data are entered in succession following period data as data to be executed within an automatic-carriage-return zone, the first space data causes the automatic carriage return operation. As a result, the second space data of the two is, against operator's will, executed to insert a space into the first column of a new line on a printing sheet of paper, whereby the beginnings of individual printed lines on the sheet of paper become uneven or indented unexpectedly.\nTo solve above problem, the existing electronic typewriters are designed not to execute the second space data if two sets of space data are entered in succession for execution within the automatic-carriage-return zone.\nThe above solution, however, has another problem which will occur when a certain number of space data are entered through a space key in order to insert the corresponding number of spaces at the beginning of a new line to open a paragraph. In other words, since space data is not executed when it is entered in the automatic-carriage-return mode, the number of space data which are executed at the beginning of a new line does not always correspond to the number of depressions of the space key. In order to open a new paragraph, therefore, it is required in the existing electronic typewriters that a desired number of space data be entered only after completion of the automatic carriage return operation. As a result, the existing electronic typewriters have a disadvantage of low typing speed."}
{"text": "String instruments include instruments such as the violin, viola, cello, double bass (sometimes called the contrabass), and harp. String instruments can be very challenging to learn and to teach, in part because mastery of string instruments requires knowledge of and experience with all of their interconnected components.\nThe violin, viola, cello, and double bass all consist of a body, a curved, hollow section made of wood where the sound resonates, and a neck, a straight piece that extends from the body with four strings stretched along it, attached to tuning pegs at the end. For example, FIG. 1 illustrates a diagram of violin and its component parts.\nPart of learning to play these instruments involves learning how to string and tune each of the musical strings. A string is made from a core, and then layers of a synthetic material or metal compound is wrapped around the core to make the string. After the string is made, a “silking” is applied. This “silking” is comprised of a colored wrapping made out of fine fibers. These fibers are wrapped at the upper and lower ends of the strings. The silking can be used to identify the brand, and to make the upper portion of the string sturdy by absorbing tension as the string is initially threaded through the hole inside of each peg and wrapped in the peg box to a desired pitch.\nFrequently, when strings are purchased in a set, there are no instructions as to how to differentiate one string from the next. An experienced musician would understand that he or she would have to separate all strings and then place them in order from thinnest to thickest in order to figure out where each string should be placed inside the peg box. However, the inexperienced musician would not know to do this without help.\nSince each peg on a fretless instrument, such as a violin, receives a specific string, a string may break because a consumer may be unfamiliar of where to place a string inside the peg box. FIG. 2 illustrates an example peg box containing the four pegs used in violins, along with the corresponding strings. Each string is a specific length and width. For example, the G string on the violin is the thickest string. The D string on a violin is wider than an A string, but an E string on the violin is the thinnest string of all four strings. If a consumer does not know how to distinguish the differences between each string, he or she may place the strings in the wrong pegs. This will cause strings to pop or break prematurely.\nTuning the instrument can also be difficult, as a string may also break when a user turns the peg past its tension point when trying to tune the string. Even when a string does not break, tuning string instruments is a problem for novice instrumentalists and to those who have difficulty with pitch recognition.\nFurther complicating the learning process is the fact that fretless instruments are unlike most other instruments because each string contains a variation of intervals and overtones. For example, on a piano keyboard, the keys are spaced in intervals consisting of either whole steps or half steps. There is no interval lower than a half step. Unlike the piano, the strings on fretless instruments contain intervals that can be played lower than half steps. The reason is because each pitch on a string is relative to where an instrumentalist places his or her fingers on the fingerboard of the instrument. Without frets, there is no guarantee that an instrumentalist will place his or her finger on the fingerboard exactly one whole step or exactly one half step from the starting pitch."}
{"text": "Field of the Invention\nThe present invention relates to a circuit configuration for protecting a power MOSFET.\nCircuit configurations for protecting power MOSFETs against overvoltage have long been known and are shown, for instance, in the Siemens Datenbuch [Siemens Databook] \"Smart-SIPMOS\", 1994/95, pages 2-7, FIG. 4. An overvoltage protection circuit configuration disclosed therein includes a Zener diode ZD2 and a diode D1.\nA MOSFET semiconductor component with a voltage spike protection element which is also known from German Published, Non-Prosecuted Patent Application DE 41 22 347 A1, includes a transistor element with a gate electrode, a source electrode, and a drain electrode. In it, a Zener diode with a cathode electrode and an anode electrode is coupled to the drain electrode of the transistor element. A horizontal MOSFET element is also presented that has a gate electrode and a drain electrode. The drain electrode is coupled to the anode electrode of the Zener diode, and the source electrode of the MOSFET element is coupled to the gate electrode of the transistor element. A back-gate electrode which is also proposed therein is coupled to the source electrode of the transistor element.\nThe wiring of the MOSFET referred to at the outset may also be achieved by using an MOS diode which is realized by an MOS transistor having a gate terminal and a drain terminal that are short-circuited, instead of using the diode D1. FIG. 2 shows one such circuit, which is described in more detail below. In particular, such a protective circuit configuration is used in high-side switches with a charge pump. The MOS diode then takes on the task of blocking the high gate voltage from the supply in the normal operating state. In the event of overvoltage, that MOS diode is operated in the conducting direction, but the on-state voltage can assume undesirably high values under some circumstances."}
{"text": "Transmissions of small utility vehicles often utilize a dog-clutch assembly operable to adjust the configuration of the transaxle. To shift gears, a user adjusts the position of a dogged shifting gear within the transmission by actuating a shifter handle. Transmissions of this type can only shift into gear when the dogs of the shifting gear and the driving gear are properly aligned. If the dogs are misaligned (a state often referred to as dead-head), the shifting gear cannot move into engagement with the driving gear.\nConventional shifting systems utilizing dog-clutches suffer from a variety of limitations and disadvantages, such as those relating to shifting gears when the transmission is dead-headed. For example, when the dogs are not aligned, the user must continue to apply force to the shifter handle until the dogs become aligned. Once the dogs become aligned, the force provided by the user causes the shifting gear to move into engagement with the driving gear, and the user can stop applying force to the shifter handle. There is a need for the unique and inventive gear-shifting apparatuses, systems and methods disclosed herein."}
{"text": "A building's structure must withstand physical forces or displacements without danger of collapse or without loss of serviceability or function. The stresses on buildings are withstood by the buildings' structures.\nBuildings five stories and less in height typically use a “bearing wall” structural system to manage dead and live load vertical forces. Vertical forces on the roof, floors, and walls of a structure are passed vertically from the roof to the walls to the foundation by evenly spreading the loads on the walls and by increasing the size and density of the framing or frame structure from upper floors progressively downward to lower floors, floor-to-floor. For ceilings and floor spans, trusses are used to support loads on the ceilings and floors and to transfer these loads to walls and columns.\nWhere vertical bearing elements are absent, for example at window and door openings, beams are used to transfer loads to columns or walls. In buildings taller than five stories, where the walls have limited capacity to support vertical loads, concrete and/or structural steel framing in the form of large beams and columns are used to support the structure.\nLateral forces (e.g., wind and seismic forces) acting on buildings are managed and transferred by bracing. A common method of constructing a braced wall line in buildings (typically 5 stories or less) is to create braced panels in the wall line using structural sheathing. A more traditional method is to use let-in diagonal bracing throughout the wall line, but this method is not viable for buildings with many openings for doors, windows, etc. The lateral forces in buildings taller than five stories are managed and transferred by heavy steel let-in bracing, or heavy steel and/or concrete panels, as well as structural core elements such as concrete or masonry stair towers and elevator hoistways.\nThere is a need for a panelized and modular system for constructing and assembling buildings without relying on concrete and/or structural steel framing, heavy steel let-in bracing, and heavy steel and/or concrete panels."}
{"text": "The present invention relates to triggering the production and introduction of shock waves into the body of a living being, such as a patient, in synchronism with the breathing of that patient.\nGerman Printed Patent Application 3,146,628 (see also U.S. Pat. No. 4,745,920 and 4,685,461) describes the comminution of concrements by means of shock waves wherein the shock waves are triggered inter alia in response to the breathing of the patient. The purpose is the following. Breathing introduces, from an overall point of view, motion and displacement of organs and parts within the body of that person. This displacement is also effective on concrements such as kidney stones to be comminuted. Hence the location of the concrements varies periodically with the breathing of the person. On the other hand, focusing of shockwaves requires a definite point and locus into which shock waves are to be focused. Therefore the motion of the concrement on the one hand, and the focusing action on the other hand, have to be synchronized, which means that in the course of oscillatory displacement of the concrement on account of breathing, one can establish the locus for the concrement during a particular phase of the breathing cycle, by triggering shockwave generation during one of the next breathing cycles in the same phase point under the assumption that the concrement will have moved back to exactly the same position it had previously during a similar phase within a breathing cycle.\nThe following are references which refer broadly to this field of art and are of an exemplary nature only. Among thim is also the U.S. Pat. No. 4,539,989 which discloses a device for comminution of concrement in connection with a coupling cushion by means of which the device producing the shock wave is acoustically coupled to the body of the person."}
{"text": "1. Field\nThe invention relates to the field of computer configuration and, more particularly, to the storage of computer configuration signals.\n2. Background Information\nComputer systems may store configuration signals in a memory. A computer system is any device including a processor capable of executing one or more instructions to generate signals. Such signals typically take the form of sequences of binary signals known as bits. Examples of computer systems are personal computers, workstation computers, server computers, hand held computers, and set top boxes to name just a few examples. Configuration signals are signals that may determine various settings for the operation of the computer system. For example, configuration signals may determine whether various input/output (I/O) ports comprised by the system are enabled, and I/O addresses for these ports. Configuration signals may determine other computer system settings as well. Such computer configuration signals are well known in the art as “set up information”. On personal computers, setup information is also often stored in a memory known as a real time clock (RTC) complementary metal oxide silicon (CMOS).\nSetup information may be applied prior to or during the booting of a program to control the computer system. Booting is the process of placing a sequence of instructions (a program) in control of various computer system resources. Resources include memory, interrupts, files, and I/O ports. An example of a program to boot is an operating system. An operating system is a program which controls various computer resources including those mentioned previously and further including typical I/O devices, such as a mouse and keyboard. Examples of operating systems are the Unix™ operating system and the Microsoft™ Windows™ operating system.\nSetup information may be read, altered, and written back to a CMOS or other memory, where it is stored using a special program called a “setup program.” The setup program may be part of the sequence of instructions comprising the computer system's power-on self test (POST) program. Often, the POST executes prior to the basic input/output system (BIOS) program of the computer system in order to initialize settings.\nThe settings determined by setup information may vary among different computer makes and models. Furthermore, the location and lengths of the bit sequences that comprise the setup information in the memory in which they are stored may vary. Accordingly, it may be difficult to create one set up program to read, alter, and write back set up) information for various makes and models of computer systems. Instead, multiple set up programs may be called for different makes and models of computer system.\nExisting setup programs typically employ a crude “textual” interface. Textual interfaces are well known in the art and may comprise an 80×25 matrix of character positions. The number, type, and position of characters in a textual interface is limited as are the number of colors in which such characters may be displayed. It is well known that such textual interfaces are more limited than modem “graphical user interfaces” (GUI) which provide individual control of the color and position on a per pixel basis on the computer system display. Typically, it is the operating system which implements a graphical user interface for the computer system. However, setup programs may execute before the operating system boots, and, therefore, the setup program may employ a less sophisticated textual interface instead of a GUI.\nThus, there is a continuing need for a setup program which may operate with various makes and models of computer systems and which may take advantage of graphical user interface features provided by an operating system."}
{"text": "Arrows have existed since before the dawn of history. The accurate placement and alignment of feathers on a shaft has long been a problem. Accordingly, numerous devices have been utilized over the centuries to apply feathers to a shaft. Some of the very recent devices have been patented.\nFor example, U.S. Pat. No. 1,896,563 to Belshaw relates to an arrow fletching machine in which horizontal movable bars operate in conjunction with an indexed wheel.\nU.S. Pat. No. 2,286,574 to Rohde relates to a fletching jig adapted to high speed production of accurately feathered arrows.\nU.S. Pat. No. 2,731,992 to Lozon relates to a fletching fixture containing a supporting structure for rotatably supporting the shaft of an arrow so that the shaft can be turned in a plurality of selective positions to have feathers applied thereto in circumferentially arranged positions.\nU.S. Pat. No. 2,742,064 to Quist relates to a fletching device for holding the shaft of an arrow in position for receiving feathers and to another device for applying a feather to the shaft.\nU.S. Pat. No. 2,836,208 to Hoyt relates to a universally adapted jig to accurately position and hold a plurality of either right or left wing feathers for either straight or spiral style fletching.\nU.S. Pat. No. 2,884,034 to Portinga relates to a fletching jig which is adjustable to accomodate the fletching of feathers upon arrows in various degrees of spiralling in either a clockwise or counterwise direction about a shaft.\nU.S. Pat. No. 2,918,097 to Thompson relates to a fletching jig for positioning the feathers on the shaft of arrows, either parallel with the shaft or an angle to the shaft.\nU.S. Pat. No. 3,024,017 to Stanton relates to providing a rigid, nonadjustable fletching jig which permits the application of an adhesive directly to the arrow shaft through the use of a slim nozzle.\nU.S. Pat. No. 3,015,483 to Martin relates to a fletching jig which allows fletches to be applied to the arrow shaft at desired angular locations."}
{"text": "1. Field of the Invention\nThe present invention pertains to a modified lignosulfonate composition; its method of preparation and its use as a dispersant, particularly in the manufacture of gypsum wallboard.\n2. Description of Related Art\nThe use of lignosulfonate as a dispersing agent in the preparation of gypsum wallboard is disclosed, for example, in U.S. Pat. No. 2,856,304. In the manufacture of gypsum products such as wallboard, lath, plaster board, sheathing or other products, calcined gypsum is formed into a slurry with a suitable amount of water and, where customary, with other additives such as paper fiber, wood fiber, starch, rosin, etc.\nIn the preparation of wallboard or similar products, in particular, the slurry then is deposited between paper liners, pressed to the desired thickness by forming rolls, allowed to set and harden, cut to desired lengths, and passed through a dryer to remove excess moisture. A portion of the water used in making the slurry combines with the calcined gypsum, as water of crystallization, in forming the final set mass of interlaced crystals, but a large portion of the water must be removed in the dryer. Obviously, the drying process is more costly as the proportion of water to be removed from the formed board is higher.\nIt is known that using lignosulfonate as a dispersing agent reduces the amount of water required to provide a flowable gypsum slurry during the deposition and forming steps. Consequently, by using lignosulfonate, the amount of water to be removed during the final drying step can be reduced, resulting in a significant economy of operation, particularly as regards lower energy costs. An ancillary benefit also often observed is higher board strength. It also is known that lignosulfonate exhibiting improved dispersing ability is prepared by base exchange of the lignosulfonate with various metals, including iron, aluminum, chromium and copper, by alkaline treatment, by oxidation and the like. See, for example U.S. Pat. Nos. 2,935,504; 3,007,910; and 3,108,008.\nUnfortunately, lignosulfonates also tend to retard the hardening or cure rate of the gypsum board, referred in the art as set retardation. While this does not present a problem with slower forming operations characteristic of the prior art, with the advent of faster processes, set retardation has matured into a significant concern. In fact, while hydrolysis and oxidation reactions tend to enhance the dispersing behavior of the lignosulfonate, such treatments tend to exacerbate set retarding characteristics.\nIn the prior art, the problem of set retardation has been dealt with primarily by adding various accelerating agents to the aqueous lignosulfonate-gypsum slurry. The prior art indicates that materials such as sodium chloride, aluminum sulfate, potassium sulfate, calcium sulfate dihydrate, and uncalcined or raw gypsum help to ameliorate set retardation.\nFinally, in U.S. Pat. No. 2,856,304 it is taught that calcining raw gypsum containing a small amount of lignosulfonate helps reduce the setting time of the calcined gypsum product.\nIn accordance with the present invention, lignosulfonate is treated in a way which enhances its capability to disperse gypsum, while at the same time, avoids imparting to the lignosulfonate an undesirable increase in its set retardation characteristics. By using the modified lignosulfonate composition of the present invention, the amount of water needed to form a calcined gypsum slurry having the necessary plasticity is considerably reduced, and the production of gypsum shapes is simultaneously accelerated."}
{"text": "1. Field of the Invention\nThe field of the invention relates to the preparation of random linear injection moldable amide-imide copolymers which process comprises reacting fully or partially acylated aromatic diamines with aromatic tricarboxylic acid anhydrides or mixtures of aromatic tricarboxylic acid anhydrides and aromatic dicarboxylic acids and to novel polytrimellitic amide-imide copolymers and to molded objects and fibers prepared from these copolymers.\n2. Background\nAmide-imide and polyamide polymers and copolymers are a relatively new class of organic compounds known for their solubility in nitrogen containing solvents when in the polyamic acid form. The major application of these amide-imides has been as wire enamels and film formers. This is illustrated in U.S. Pat. Nos. 3,852,106 (1974), 3,817,942 (1974), 3,661,832 (1972), 3,454,890 (1970) and 3,347,942 (1967).\nPolyimide and polyamide-imide polymers have also been found useful for molding applications as shown in U.S. Pat. Nos. 4,016,140 (1977), 3,654,227 (1972) and 3,573,260 (1971).\nThe general object of this invention is to provide injection moldable amorphous linear amide-imide copolymers. A more specific object of this invention is to provide a novel process for the manufacture of injection moldable amide-imide and amide copolymers by reacting fully or partially acylated aromatic diamines with aromatic tricarboxylic acid anhydrides or mixtures of aromatic dicarboxylic acids and aromatic tricarboxylic acid anhydrides. Another object is to provide novel polyamide-imide copolymers suitable for use as an engineering plastic particularly for and in injection molding. Other objects appear hereinafter.\nWe have discovered a novel melt condensation process in which fully or partially acylated aromatic diamines are reacted with aromatic tricarboxylic anhydrides or mixtures of aromatic tricarboxylic anhydrides with aromatic dicarboxylic acids to yield engineering plastics suitable for injection molding which feature very high tensile strength and heat distortion temperatures. Our novel process for the preparation of random linear injection moldable amide-imide and amide copolymers comprises reacting fully or partially acylated aromatic diamines with aromatic tricarboxylic acid anhydrides or mixtures of aromatic tricarboxylic acid anhydrides with aromatic dicarboxylic acids in a molar ratio of about 9:1 to about 1:9 at a temperature of about 450.degree. to 750.degree. F.\nIn the prior art, melt reaction of tricarboxylic acid anhydride compounds with aromatic diamines have produced products which are not suitable for injection molding application. The reason for this is not known, but it is specified that various side reactions occur. It has now been discovered that when fully or partially acylated diamines are reacted, injection molding grade polyamide-imide copolymers are produced. In our process we usually acylate more than half of the diamines utilized in the reaction. The preferred acylation is about 70 to 100 percent.\nEvidence for linearity for our novel copolymer is demonstrated by the solubility of the polymer. Polymers produced from tricarboxylic acid anhydride compounds such as trimellitic acid anhydride and aromatic diamines via various melt polymerization methods show no solubility for products having inherent viscosity in excess of 0.5. The copolymer produced according to the novel process utilizing partially or fully acylated diamines is essentially soluble after curing with inherent viscosities in the range of 0.6 to 3.0. For the purpose of this invention, inherent viscosity is measured at 25.degree. C. and 0.5% w/v in 100% sulfuric acid or N-methylpyrrolidone.\nThe novel injection moldable amorphous random linear polyamide-imide copolymers of this invention comprise units of ##STR1## R comprises R.sub.1 and R.sub.2, R.sub.1 and R.sub.2 are divalent aromatic hydrocarbon radicals of from 6 to about 20 carbon atoms or two divalent hydrocarbon radicals of from 6 to 20 carbon atoms joined directly or by stable linkages selected from the group consisting of --O--, methylene, --CO--, --SO.sub.2 --, and --S-- radicals and wherein said R.sub.1 and R.sub.2 containing units run from about 10 mole percent R.sub.1 containing unit and about 90 mole percent R.sub.2 containing unit to about 90 mole percent R.sub.1 containing unit and about 10 mole percent R.sub.2 containing unit.\nThe novel injection moldable random linear copolymer may comprise structural Units A and B and also include Unit C of the following formula: ##STR2## wherein X is a divalent aromatic radical usually a divalent benzene radical and R.sub.3 comprises both R.sub.1 and R.sub.2 as defined above or is equal to R.sub.1. Furthermore, if structure C is present R of structural Units A and B can be equal to R.sub.1 or comprise both R.sub.1 and R.sub.2 as set forth above.\nIn the foregoing structural units Z is a trivalent aromatic radical. Z may be a trivalent radical of benzene, naphthalene, biphenyl, diphenyl ether, diphenyl sulfide, diphenyl sulfone, ditolyl ether, and the like.\nUseful aromatic tricarboxylic acid anhydrides which contribute the trivalent radical moiety of Z include those compounds containing at least one pair of carboxyl groups in the ortho position with respect to each other or otherwise situated in a fashion which permits the formation of an anhydride structure, one other carboxyl group and from 9 to 21 carbon atoms. Within these limits, these compounds may contain one or more benzenoid rings such as, for instance, trimellitic anhydride and its isomers and multi-ring compounds such as the 1,8-anhydride of 1,3,8-tricarboxylnaphthalene. Usually these compounds contain up to three benzenoid rings.\nThe aromatic tricarboxylic acid anhydride used in the novel process to form the polyamide-imide polymers of this invention is of the formula: ##STR3## where Z is a trivalent aromatic radical defined as set forth hereinabove. The following aromatic tricarboxylic anhydrides are preferred: trimellitic acid anhydride; 2,3,6-naphthalene tricarboxylic anhydride; 1,5,6-naphthalene tricarboxylic anhydride, and the like; 2,6-dichloronaphthalene-4,5,7-tricarboxylic anhydride, and the like. One of the preferred aromatic tricarboxylic anhydrides is trimellitic anhydride since this compound is readily available and forms polymers having excellent physical properties of tensile strength and elongation and is resistant to high temperatures.\nSuitable fully or partially acylated aromatic diamines useful in applicant's process include para- and meta-phenylenediamine, oxybis (aniline), thiobis (aniline), sulfonylbis (aniline), diaminobenzophenone, methylenebis (aniline), benzidine, 1,5-diaminonaphthalene, oxybis (2-methylaniline), thiobis (2-methylaniline), and the like. Examples of other useful aromatic primary diamines are set out in U.S. Pat. No. 3,494,890 (1970) and U.S. Pat. No. 4,016,140 (1977) both incorporated herein by reference.\nUseful aromatic dicarboxylic acids include isophthalic acid and terephthalic acid. In applicant's process further preparation of injection moldable amide-imide and amide copolymers process can be conducted without utilizing a solvent or fluidizing agent though it is preferred to use agents such as N-methylpyrrolidone, dimethyl-acetamide, or acetic acid for the initial mixing of reactants. In general, since these polymers are linear, they may be easily cured in the melt using a twin screw extruder as the finishing reactor instead of a solid state polymerization. However, in some instances, it may be helpful to solid state polymerize the copolymers. The term \"solid state polymerization\" refers to chain extension of polymer particles under conditions where the polymer particles contain their solid form and do not become a fluid mass.\nThe solid state polymerizing can be carried out below the melting point of the copolymer and can be conducted in several ways. However, all the techniques require heating the ground or pelletized copolymer below the copolymer melting point, generally of about 400.degree. to 600.degree. F. while either sparging with an inert gas such as nitrogen or air or operating under vacuum.\nInjection molding of the novel copolymer is accomplished by injecting the copolymer into a mold maintained at a temperature of about 350.degree.-500.degree. F. In this process a 0.1-2.0 minutes cycle is used with a barrel temperature of about 500.degree. F. to 700.degree. F. The injection molding conditions are given in Table I.\nTABLE I ______________________________________ Mold Temperature 350-500.degree. F. Injection Pressure 2,000-40,000 psi and held for 0.5-20 seconds Back Pressure 0-400 psi Cycle Time 6-120 seconds Extruder: Nozzle Temperature 500.degree. F. to 700.degree. F. Barrels: Front heated to 500.degree. F. to 700.degree. F. Screw: 10-200 revolutions/minute ______________________________________\nThe mechanical properties of the polymers prepared in the Examples are given in Tables II, III, IV, V, VI and VII.\nIn applicant's process the acylated aromatic diamines need not be isolated or purified prior to their further reaction with the tricarboxylic acid anhydride compound or mixture of the tricarboxylic acid anhydride with dicarboxylic acid. Therefore, one can react one to two moles of acetic anhydride or acid or propionic anhydride or acid or any other C.sub.2 through C.sub.8 containing aliphatic anhydride or acid and one mole of the appropriate aromatic diamine or diamine mixture and use the resulting diacylated diamine solution in acetic acid or propionic acid to react with the tricarboxylic anhydride compound, or mixtures of the tricarboxylic anhydride compound with dicarboxylic acid.\nIn most cases, linear high molecular weight polyamide-imide polymers result after melt and/or solid state polymerization."}
{"text": "Technical Field\nThe present disclosure relates generally to the field of reposable or reusable surgical instruments. In particular, the disclosure relates to instruments having separable and replaceable components to provide clean, sterile or refurbished surfaces in each instance of use.\nBackground of Related Art\nOne type of surgical device is a linear clamping, cutting and stapling device. Such a device may be employed in a surgical procedure to resect a cancerous or anomalous tissue from a gastro-intestinal tract. Conventional linear clamping, cutting and stapling instruments include a pistol grip-styled structure having an elongated shaft and distal portion. The distal portion includes a pair of scissors-styled gripping elements, which clamp the open ends of the colon closed. In this device, one of the two scissors-styled gripping elements, such as the anvil portion, moves or pivots relative to the overall structure, whereas the other gripping element remains fixed relative to the overall structure. The actuation of this scissoring device (the pivoting of the anvil portion) is controlled by a grip trigger maintained in the handle.\nIn addition to the scissoring device, the distal portion also includes a stapling mechanism. The fixed gripping element of the scissoring mechanism includes a staple cartridge receiving region and a mechanism for driving the staples up through the clamped end of the tissue against the anvil portion, thereby sealing the previously opened end. The scissoring elements may be integrally formed with the shaft or may be detachable such that various scissoring and stapling elements may be interchangeable.\nA number of surgical device manufacturers have developed product lines with proprietary drive systems for operating and/or manipulating such surgical device. In many instances, such surgical device further includes a handle assembly, which is reusable, and a disposable end effector or the like that is selectively connected to the handle assembly prior to use and then disconnected from the end effector following use in order to be disposed of or in some instances sterilized for re-use.\nSurgical devices that are reposable, or reusable for multiple procedures, reduce the instrumentation costs per procedure. Providing a reusable surgical device, however, presents various challenges. For example, the complexity of a surgical device tends to result in fairly labor intensive cleaning procedures to prepare the surgical device for subsequent use. Improper cleaning may result in dangerous contamination being introduced into the surgical site. Also, some reusable surgical devices have removable and replaceable components to provide clean surfaces for each use. Many of these surgical devices require arduous disassembly and reassembly procedures that require extensive training, and may discourage use of the surgical device."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for fabricating a Alxe2x80x94Si alloy packaging material for semiconductor device, and in more particular to a method for fabricating an Alxe2x80x94Si alloy packaging material which is capable of increasing Si content by mixing Si powders with Alxe2x80x94Si powders and pressurizing-forming the mixture or by pressurizing-filling Si powders or a preforming body of Si powders with Alxe2x80x94Si alloy melt.\n2. Description of the Prior Art\nA packaging material for semiconductor device means a material fabricated as a box-shaped body in order to protect a semiconductor device disposed therein. In general, a package is largely divided into a homogeneous package and a heterogeneous package.\nFIG. 1a illustrates a homogeneous package consisting of a base 1, a side-wall 2 and a lid 3. Herein, a feedthrough 4 is arranged between the base and the side-wall 2, and semiconductor substance 5 is disposed in the package. FIG. 1b illustrates a heterogeneous package fabricated by joining the base 1, the side-walls 2 and the lid 3, herein, feedthrough 4 is formed on the sidewalls 2.\nThere are required characteristics for a packaging material for semiconductor device.\nFirst, a packaging material has to have a thermal expansive coefficient similar to that of a semiconductor device. In more detail, when the semiconductor device is operated, heat is generated, and accordingly a temperature of the semiconductor and the package rises, herein, if they have different thermal expansive coefficients, the semiconductor device may be dissected from the package.\nIn addition, it is preferable for the packaging material to have a high heat transfer modulus. When the packaging material has a high heat conductivity, heat generated in the operation of the semiconductor device can be quickly emitted to the outside, and accordingly it is advantageous in maintaining performance of the semiconductor device.\nThe lighter the packaging material, the more it is advantageous. By using the packaging material having a low density, a weight of electronic appliance can be reduced, and accordingly mobility of electronic appliance can be improved.\nThe packaging material is required to have a good processability. In more detail, in order to fabricate an intricate package or a package required to have a precise measure tolerance, the packaging material having a good processability has to be used.\nIn addition, the packaging material is required to have a good plating characteristic. If a surface of the packaging material is plated with nickel, copper, silver and gold, etc. by an electroplating method and the joining between the packaging and the plating layer is strong, a life-span of the packaging can be improved.\nBesides the above-mentioned characteristics, it is preferable to use a packaging material having a good bonding characteristic. In joining of the base, the side-walls and the lid, by using a material having a good bonding characteristic, a life-span and bonding characteristics of the packaging can be improved. Herein, the base, the side-walls and the lid are fabricated with the same material or different materials, and welding or soldering or adhesion agent can be used to join the construction parts.\nLastly, it is required for the packaging material to have a simple fabrication process and low production cost.\nVarious packaging materials having all those characteristics have been developed, among them typical packaging materials will be described with reference to following table 1.\nThe above-mentioned packaging materials respectively have merits and demerits. For example, kovar has an appropriate thermal expansive coefficient, however, its heat conductivity is low and a density is high. Titanium has a good thermal expansive coefficient and a low density, however, it has low heat conductivity and is expensive. Aluminum has an appropriate heat conductivity and a low density for a packaging material, however, because it has a high thermal expansive coefficient, its usage for a packaging material is restricted. Beryillia and diamond respectively have a low thermal expansive coefficient, a high high heat conductivity and a low density, however, they are excessively expensive.\nIn the meantime, a metal matrix composite fabricated by adding a stiffening agent such as SiC, B4N3 and Al2O3, etc. to a metal material such as aluminum, etc. shows good characteristics as a packaging material in several aspects such as a thermal expansive coefficient, a heat conductivity, a density and a price, etc.\nFIG. 2 illustrates a pressureless infiltration method as one of typical methods for fabricating Alxe2x80x94SiC one of metal matrix composites. In the pressureless infiltration method, in order to make metal melt 8 easily penetrate into a preforming SiC body 7, a certain processes are required. In more detail, in order to improve a wetability of ceramic to the metal melt, a coating layer such as oxide, etc. is formed onto the ceramic surface or a chemical reaction between the melt and the ceramic surface is induced by using a special bonding agent. Afterward, the melt can easily penetrate into the preforming ceramic body. However, the above-described method is intricate. In addition, ceramic has bad processing characteristics.\nIn the meantime, in consideration of a thermal expansive coefficient, a heat conductivity and a density, etc., Si-30 wt % Al alloy is appropriate for a packaging material. In addition, its material cost, processability and bondingability are good, its application range will be gradually expanded.\nRecently, a method for fabricating Alxe2x80x94Si alloy having the content of Si in the range of 50xcx9c70% by total weight of Alxe2x80x94Si alloy with a low cost has been developed in the U.K. FIG. 3 illustrates the conventional method for fabricating an Alxe2x80x94Si alloy packaging material. In the method, by fabricating preforming alloy Alxe2x80x94Si melt 9, spraying it as liquid 10 having a 50 micronxcx9c200 micron size by using high pressure gas such as high pressure nitrogen and argon, etc., the alloy having a billet shape, etc. is formed onto a substrate moving rotatively and horizontally. After cuffing the formed Alxe2x80x94Si alloy body 11 so as to have a certain width, a hydrostatic press process (HIP, process for eliminating internal air holes), a mechanical process, a plating process, a semiconductor device adhesion process, and a packaging sealing process, etc. are sequentially performed. FIG. 4 is an enlarged photograph of the Al-70Si alloy fabricated by a spray forming method of OSPREY company in the U.K. taken with an electronic microscope ZEISS Axioskip (Germany).\nHowever, the conventional method has following demerits. In the conventional method, a temperature of the melt has to be not less than 1000xc2x0 C., in cutting of the formed body and fabricating it as a certain package box, material loss occurs, and accordingly a fabrication cost is increased. In more detail, in fabricating of Alxe2x80x94Si packaging material by the conventional spray forming method, because the more the Si content, the higher a temperature of Alxe2x80x94Si melt has to be risen, there is lots of energy loss. In addition, because the spray-fabricated Alxe2x80x94Si alloy body has to be cut again to have a request size, lots of material and time loss occur.\nIn order to solve the above-mentioned problems, it is an object of the present invention to provide a method for fabricating an Alxe2x80x94Si packaging material which is capable of simplifying fabrication processes and lowering a production cost in comparison with the conventional Alxe2x80x94Si fabrication method by mixing Si powders with Alxe2x80x94Si alloy powders and filling the mixture as a request size from beginning.\nIt is another object of the present invention to provide a method for fabricating an Alxe2x80x94Si packaging material which is capable of fabricating an Alxe2x80x94Si alloy material having lots of Si content by easily adjusting the Si content in fabrication.\nIn addition, it is yet another object of the present invention to provide a method for fabricating an Alxe2x80x94Si packaging material which is capable of reducing internal air holes and performing full filling without using outer big stress by mixing-filling Alxe2x80x94Si alloy powders having a high temperature flowability with Si powders not having plastic deformation characteristic.\nIn addition, it is still another object of the present invention to provide a method for fabricating an Alxe2x80x94Si packaging material having different Si content in the thickness direction by differentiating a mixing proportion of Si of each or both Si powders and Alxe2x80x94Si alloy powders."}
{"text": "Laser photocoagulation of the retina has been practiced for more than 30 years. Its therapeutic success is undisputed, being founded on many studies for different retinal diseases. The uses of the laser are both therapeutic (for example diabetes, thromboses of the eye, age-related macula degeneration (AMD)) and preventive (for example retinopexy). For photocoagulation, according to the type of disease, between a few single expositions in the macula and up to 3000 foci are applied in the case of pan-retinal photocoagulation. According to the spot size, 100-500 mW laser power is applied within 50-300 ms per focus. Mainly lasers in the green spectral range (Ar-ion lasers with 514 nm or frequency-doubled Nd-lasers at 532 nm) are used, but also lasers and laser diodes in the near IR-range.\nThe only dosage of laser photocoagulation employed so far consists in subsequently checking the ophthalmologic appearance of the coagulation site at the fundus. That the retina turns grey or white in the process, shows an irreversible thermal necrosis of the neuronal retina that can reach, depending on the intensity and extent of the coagulation, the entire retina from the retinal pigment epithelium (RPE) via the photoreceptors to the nerve fiber layer and can also include necroses of all uvulas. The large spatial extent of the coagulation effects results from the thermal conduction from the melanine-containing absorbing layers into the neighboring tissue layers.\nThe pan-retinal laser coagulation of the retina is the most common use of the laser in ophthalmology. Here the inner layers of the peripheral retina are to be destroyed thermally in 3-10 sessions with up to 3000 laser foci of different size so as to prevent the unchecked vessel growth and the blindness connected thereto at a later point. With most patients the laser treatment is extremely painful. Only a retrobulbar injection can avoid pain. This however entails that the mobility of the bulbus that is necessary for carrying out the photocoagulation in the periphery is switched off: in the case of over-coagulations and in particular in the case of repeat treatments, additional thermal damage of ganglion cell layers is to be feared that can lead to extensive defects of the field of vision.\nSince the absorbing granula vary considerably in terms of their local and spatial density it is not surprising that the histological results after laser coagulation can vary considerably even in the case of identical exposition parameters. The extent of the damage essentially is a function of the extent of the laser-induced rise in temperature. In practice, it cannot be predicted due to different pigmentation and thus absorption of the retina both inter- and also intraindividual.\nAutomatic, temperature-controlled online dosimetry for laser treatment with minimal invasive damage is a goal to be desired that cannot be achieved through the application method that is conventional at the moment for ophthalmoscopic examination.\nPhysics offers different methods for measuring temperatures, however almost all of them are practically unsuitable for measuring the ocular fundus.\nInvasive measurement methods as for example thermal probes or dyes that fluoresce as a function of the temperature are too annoying—among others on account of side effects—and/or too imprecise. Due to the absorption of the thermal radiation in the eye, thermal imaging cameras cannot be used.\nMethods that are based on auto-fluorescence seem to be suitable, as is for example taught by DE 102 40 109 A1, but a uniform distribution of the chromophores that does not exist in practice is a precondition here.\nThe analysis of temperature-dependent, thermo-mechanical expansion of an absorber and the pressure wave emitted therewith after the application of a short laser pulse has been described in Sigrist M. W., “Laser Generation of Acoustic Waves in Liquids and Gases”, Journal of Applied Physics 60(7):R83-R121, 1986. On this basis the optoacoustic temperature measurement on the retina was developed as it is illustrated in DE 101 35 944 C2. Additional, repetitive irradiation with short laser pulses produces pressure transients whose amplitude can be recorded with an ultrasound sensor (for example piezo element) that is integrated into the contact lens required for laser treatment anyway. The instantaneous increase in temperature can be determined from the amplitude. In the process, the dependency of the temperature on the choroid perfusion and the light absorption and thus the absolute necessity of an online dosimetry based on the temperature could be illustrated.\nThe method of DE 101 35 944 C2 was previously used for temperature measurements in Transpupillary Thermotherapy (TTT) and the Selective Retina Therapy (SRT). In the case of SRT, the treatment pulses themselves can be used for determining the temperature. WO 2005/007002 A1 further describes strategies that use the optoacoustic signal for controlling the therapeutic laser. WO 2005/007002 A1 however assumes that microscopic bubbles will form sooner or later due to the laser impact. Such bubbles significantly change the behavior of the pressure transients and thus serve to identify the damage threshold in the vicinity of which the laser should operate. This is then realized by a suitable feedback.\nIn laser photocoagulation, only maximum temperatures of 40-80° C. are realized for treatment of the ocular fundus. Bubble formation cannot set in below 100° C., so that this cannot be any option for laser control.\nThe aim of photocoagulation is the thermal denaturation of proteins and tissue. It is in particular the dependency of the extent and depth of damage of coagulations of the retina that has been well researched experimentally and theoretically using different laser parameters (e.g. Birngruber R, Hillenkamp F, Gabel V P., “Experimental studies of laser thermal retinal injury”, Health Phys 44(5):519-531, 1983 or Birngruber R, Hillenkamp F, Gabel V P., “Theoretical investigations of laser thermal retinal injury”, Health Phys 48(6):781-796, 1985). The findings are that the damage to the tissue is both a function of the duration of the laser irradiation and also directly—and particularly critically—of the temperature increase caused during this period. Here the damage integral Ω describes a certain change that is a function of damage criteria and tissue and that is influenced by the temperature curve T(t) over the total duration of the temperature increase ts.\n      Ω    ⁡          (              t        s            )        =      A    ·                  ∫        0        ts            ⁢                          ⁢                        ⅆ          t                ·                  T          ⁡                      (            t            )                          ·                  ⅇ                      -                                          Δ                ⁢                                                                  ⁢                E                                            k                ·                                  T                  ⁡                                      (                    t                    )                                                                                          \nThe activation energy ΔE and the frequency factor A can be determined experimentally, in that the threshold for thermal damage for different temperature increases and exposition times are determined, k designating the Boltzmann constant. The constants differ for many tissues. Ω is influenced exponentially by the temperature and approximately linearly by the time. This means that the effects of an excessive temperature can be far more serious than that of an irradiation period that is too long. To describe a stronger denaturation, Ω>>1 is selected (e.g. Ω=100), if the value clearly stays below Ω=1, Ω<<1, no thermal changes are to be expected."}
{"text": "Various trailers, such as, by way of non-limiting example, flat-bed trailers, utility trailers, travel trailers, are typically coupled via a trailer tongue to a towing hitch associated with a towing vehicle. The towing hitch may include a ball hitch that is received in a ball receiver associated with the trailer tongue. Additionally, the trailer tongue may include one or more secondary connecting components, such as, by way of non-limiting example, trailer chains and/or an electrical plug.\nWhile the trailer is being towed, the secondary connecting components may couple and/or connect the trailer to a portion of the towing hitch. By way of non-limiting example, an attachment portion of each of the one or more trailer chains may be received by a portion of a towing hitch. The attachment portion may be a hook or a clasp that is configured to be received by an aperture disposed on the portion of the towing hitch. The one or more trailer chains may act as a safety feature of the trailer. The one or more trailer chains may act to secure the trailer to the towing vehicle if the ball hitch is dislodged from the ball receiver.\nThe towing vehicle may include a wiring harness. The wiring harness may be configured to receive an electrical plug associated with the trailer. The wiring harness and the electrical plug may cooperatively operate to synchronize various lights on the towing vehicle with various lights on the trailer. For example, the wiring harness and electrical plug may cooperate to synchronize turn indicator lights and break lights on the towing vehicle with turn indicator lights and break lights on the trailer.\nWhen the secondary connecting components of the trailer are not in use, the secondary connecting components may drag or lay on the ground. For example, one or more trailer chains may drag on the ground when the trailer is being towed and the one or more trailer chains are not coupled to the towing vehicle. Additionally, the one or more trailer chains may lay on the ground when the trailer is not being towed. Similarly, the electrical plug may drag on the ground when the trailer is being towed and the electrical plug is not connected to the wiring harness and/or lay on the ground when the trailer is not being towed. The secondary connecting components may become damaged as a result of exposure to dirt, moisture, or friction from laying on the ground and/or being dragged on the ground. Additionally, the secondary connecting components may have to be lifted off of the ground in order to couple and/or connect them to the towing vehicle. Accordingly, a system for storing the secondary connecting components while the secondary connecting components are not in use may be desirable."}
{"text": "The present invention relates to a roller valve lifter having a roller at one end thereof that cooperates with a lobe of a camshaft in an internal combustion engine. More specifically, the invention relates to improving the lubrication of the valve lifter and preventing rotation of the lifter.\nConventional camshaft or xe2x80x9ccamxe2x80x9d, internal combustion engines typically utilize valve lifters, push rods, and valve springs along with rocker arms to open and close the valves of the engine to allow air and fuel to enter and exhaust to exit the cylinders of the engine during combustion. These components are collectively referred to as the xe2x80x9cvalve train.xe2x80x9d\nIn conventional cam engines as opposed to those of over-head design, a valve lifter with a pushrod rides on the lobes of the camshaft which is rotated by the crankshaft. As the lifter reciprocates up and down, the push rod seated in the lifter also reciprocates and communicates this up and down motion via a rocker arm to either an intake or exhaust valve. A high tension spring ranging from approximately 200 to 1000 ftxc2x7lbs, surrounds the stem of the valve and when the spring is compressed, the valve is pushed into the cylinder.\nDuring the up stroke of the piston in the cylinder, the intake valve opens to allow fuel and air to enter the combustion chamber. Somewhere near the very top of the up stroke, both the intake and the exhaust valves close and the spark plug creates a spark to ignite the air-fuel mixture which is under compression by the piston. This results in a high temperature explosion which forces the piston downward, called the xe2x80x9cpower stroke,xe2x80x9d thereby translating this movement via a connection rod to rotate the crankshaft which, in turn, translates this angular motion to the wheels of the vehicle via a set of gears. Near the bottom of the compression stroke, the exhaust valve opens to expel the burnt fuel mixture out of the cylinder. After the piston changes directions and begins the up stroke, the exhaust valve continues to remain open thereby forcing any remaining the spent gases out of the cylinder. However, during this same time, the intake valve begins to open to recharge the cylinder with fuel. It is not until the piston has started to travel upward that the exhaust valve closes. Thus, at various times during the compression cycle, both the intake and exhaust valves will be open and closed at the same time. The timing of the opening and closing of the valves is controlled by the physical design of the oval shaped lobes on the camshaft. As the valve lifter is pushed upward by the lobe of the camshaft, the valve lifter pushes the pushrod up which drives the rocker arm downward, causing the valve to open. Likewise, as the lifter and pushrod travel downward, the rocker arm raises and the valve closes due to the biasing action of the valve spring.\nIn high speed engines, measured as revolutions per minute (RPM), the valve train components are under extreme stress and high temperatures. To increase engine performance and decrease component wear which may eventually lead to failure, various valve lifter configurations have been designed. Solid and hydraulic valve lifters are the most common designs used in conventional cam engines. Hydraulic lifters are typically used in relatively low RPM engines, up to 6,500, whereas solid valve lifter designs are preferred in high RPM applications such as racing and high performance applications. Conventional hydraulic and solid lifters have a flat surface that is fixed or integral with the body of the lifter and is adapted to engage and ride on the lobes of the camshaft. The engagement between the fixed surface of the lifter body and the camshaft lobe creates high frictional forces causing the surfaces of the lobes to wear. Therefore, the higher the RPM of the engine, the greater the wear and the likelihood of material being removed. As material is removed from the surface of the lobe, the timing of the opening and closing of the valve also changes. This change in timing may hamper engine performance such as by allowing excess air-fuel mixture to enter the cylinder causing a rich condition. Conversely, improper timing may permit the air-fuel mixture to escape through the exhaust valve which results in a lean condition. Either of these conditions will affect cylinder pressure and decrease performance and may cause misfiring of the cylinder and engine damage. Furthermore, if this improper timing allows a valve to remain open when the piston is near the top of the compression stroke, the piston will strike the valve resulting in bent pushrods and valves, broken valve springs and lifters and will eventually lead to catastrophic engine failure.\nTo decrease lobe wear in high performance engines, a roller has been added to the body of the valve lifter for riding on the cam. The roller allows the use of a camshaft with lobes of steeper ramp angles to provide faster valve opening and closing for accommodating high RPM engines. The roller engagement between the roller and rotating cam lobe reduces the frictional forces generated therebetween. Not only does the presence of the roller decrease cam lobe and valve lifter wear, it also provides smoother transitions as the roller travels over the peak of the lobe thereby decreasing valve train noise. Likewise, various bearing and sleeve configurations have been utilized to decrease friction and wear of the shaft rotatably mounting the roller to the valve lifter. For high performance engines, needle bearings have replaced solid rollers, bushing and conventional ball bearings to decrease wear and more evening spread the load over the surface of the shaft. However, these bearings and bushings also rely upon oil to function properly.\nAlthough the addition of the roller increases camshaft and valve train life, overall roller wear is a function of engine speed (RPM). High performance engines such as those used in drag racing applications produce extremely high engine speeds (6,000 to 13,000 RPM) over a short duration of time (i.e. less than 5 to 12 seconds). Conversely, stockcar racing engines produce relatively high engine speeds of typically 5,000 to 8,000 RPM and under racing conditions, maintain those speeds for long periods of time (2 to 3 hours). At these high engine speeds, it becomes difficult to provide oil to the valve lifter, roller and bearing assembly as well as adequate lubrication of the camshaft.\nFrom the ground up, a typical engine is configured with an oil pan for holding oil and an oil pump which feeds the oil to various locations in the engine. Above the oil pan sits the engine block and the crankshaft, such that a portion of the crank rotates in the oil. In a typical xe2x80x9cVxe2x80x9d-style engine, that is, one having cylinders at an angle to the left and right sides of the block in a xe2x80x9cVxe2x80x9d pattern with the crankshaft positioned at the apex of the xe2x80x9cVxe2x80x9d, the camshaft is typically located directly above and in parallel with the crank. In straight cylinder configuration engines wherein all cylinders are aligned in a row, the crankshaft and cylinders are located in the same plane and camshaft is positioned to one side so not to interfere with the travel of the connecting rods.\nThe valve lifters, in an xe2x80x9cVxe2x80x9d style engine, are located in a lifter galley. The lifters are lubricated by oil in the engine block and receive direct lubrication from a transverse oil passageway in the engine block that intersect the bores in which the valve lifters are positioned and indirectly from oil that is sprayed into the lifter galley from the rotation of the crankshaft and connecting rods. Various methods have been employed to increase the lubrication of the valve lifters and camshaft.\nOne method used to increase the movement of oil to the valve lifters and camshaft is the addition of small holes to the crankshaft and the dynamic balance weights of the crank. These holes, or oil squirters, pickup oil from the pan and any oil on the surface of the crank and throw the oil to the camshaft and valve lifter as the crankshaft and rotates. This method is also employed in engines having steel connecting rods to lubricate the cylinder wall by placing a through-hole on the end that connects to the piston and to the lifters by adding a squirter to the xe2x80x9cbig endxe2x80x9d or end that connects to the crankshaft. However, the machining of the squirter reduces the strength of the connecting and have been found to severely weaken aluminum connecting rods used in high performance, high RPM engines.\nAnother method of directing oil to the lifters and camshaft involves adding separate oil feed lines to the lifter galley. This is accomplished by drilling a feed hole into an oil passageway of the engine block to tap the oil pressurized by the oil pump and adding metal tubing to direct the oil to the desired location such as above the camshaft. However, adding components to the internals of engine is not always practical due to the limited amount of space. Furthermore, these added components may also fail and create shrapnel that will be run through the engine which can damage precision surfaces such as on the camshaft, crankshaft, pistons, etc.\nTo increase the movement of oil in the common transverse oil passageway and lifter bores, the valve lifter body has been modified. One modification includes adding a channel through the body of the lifter to increase the amount of flow of oil from one passageway to the next lifter bore. Another method of facilitating the flow of oil in the common passageway while increasing lubrication to the lifter is by adding an annular groove to the body of valve lifter. As the valve lifter reciprocates in the bore, the oil trapped between the space created by the annular groove and the bore is deposited on the walls of the bore.\nWith all of these methods, the higher the RPM, the greater the oiling of the valve lifter; however, at low engine speeds such as during idling, start-up, stop-and-go driving conditions, and gear shifting create inadequate lubrication conditions. Not only are these types of driving conditions prevalent on race day, but also seen during every day driving. Therefore, a method is needed to provide adequate lubrication to the roller and the bearing assembly thereof to reduce wear, maximize engine performance and avoid valve train component failure.\nAnother problem associated with the use of solid valve lifters with rollers in high RPM engines, is the rotation of the lifter as it reciprocates in the lifter bore of the engine. At high RPM the valve lifter has a tendency to rotate so that its axis of rotation becomes skewed or out of parallel alignment with that of the camshaft and lobes thereof. Also, the use of steep angled camshaft lobes require extremely high valve spring pressures. Any misalignment of the roller with the engaging surface of the camshaft lobe may lead to catastrophic failure of the roller causing significant damage to the camshaft and bent pushrods and valves and broken rocker arms and valve springs. Also, rotation of the lifter in the bore may prevent the oil pressure feed receiving area or groove of the valve lifter from intersecting and the common transverse oil passageway of the engine block that feeds oil to the valve lifters.\nTo prevent rotation in the bore, link bars are commonly used to tie the bodies of two lifters together, typically the exhaust and intake of one cylinder. These link bars may be permanently attached to the lifters or removable such as shown in U.S. Pat. No. 4,809,651. Although these prior link bars prevent rotation, they also add components and weight to the lifter assembly. Furthermore, the attachment point of the link bar to the body also wears due to the repetitive motion and may eventually fail. Furthermore, in high revolutions engines, these link bars on the valve lifters are constantly fighting rotation and under repetitive forces. Thus, in applications requiring high engine speeds over long durations of time, the link bar and the attachment devices may fatigue creating unnatural movement of the lifter which will damage the valve train.\nAnother method used to prevent rotation of the lifter is by adding a xe2x80x9cUxe2x80x9d shaped member in which the legs of the xe2x80x9cUxe2x80x9d are inserted into two adjacent lifter bores as illustrated in U.S. Pat. No. 5,022,356. The legs of this anti-rotation member are smaller than the diameter of the lifter bore and longer than the bore length. Once inserted in the lifter bore, the member is push to the front or rear of the bore and, thus, the member makes contact with the entire length of lifter bore on each end side of the member leg. The member is prevented from falling through the bores by a cross-member that connects the two legs. Also, a foot is added at the end of the member to prevent the member from exiting the lifter as the lifter travels upward. The valve lifter must also be modified to be used in conjunction with this member. The portion of the valve lifter which engages the member must be machined flat. Although this member and lifter assembly prevents rotation without adding components to the valve lifter body, the member presents other problems. The member edges are in contact with the full length of the lifter bore and the long flat of the valve lifter engages the member. Thus, as the lifter reciprocates, the large area of contact between the member and the lifter creates friction thereby requiring additional lubrication to prevent excessive wear and heat. Furthermore, the edges of the member may eventually wear into the lifter bore thereby removing material which is run through the engine. Also, the feet of the member extend through the lifter bore positioning themselves near the camshaft and the roller of the lifter. The height of the feet are, therefore, critical to prevent the lobes of the cam from making contact with them. In high performance engines, a specific cam design is used to create precise opening and closing of the valves for that particular engine configuration. Thus, if an engine is retrofitted with a different camshaft, the feet of the member may also have to be ground to allow clearance by the cam lobes. Therefore, an anti-rotation device which prevents rotation of the lifter but does not add weight and/or components to the valve lifter or those that may interfere with the cam lobes and does not create excess friction and heat is needed for these high performance engines.\nIn accordance with the presence invention, a valve lifter apparatus is provided including a body with a roller member at one end thereof for riding on one of the camshaft lobes. The body is provided with a predetermined flow path which direct lubrication in a well-defined manner directly to be end of the body at which the roller member is located. In this manner, lubrication is directed in a predetermined manner to the place it is needed most, i.e. the roller, rather than simply relying on the general undirected travel of the oil fed to the lifter bore.\nIn another aspect of the invention, a valve lifter assembly is provided including a lifter body which reciprocates in a bore in the engine block. A portion of body of the valve lifter has a flat exterior surface and the assembly includes an anti-rotation member including at least one short portion thereof that extends into the lifter bore adjacent the flat of the valve lifter body to prevent rotation thereof in the bore. As the length of the flat is much greater than the length of the member portion, the flat surface will only have a short section thereof that is in contact with the short member portion at any time during the reciprocation of the valve lifter body. This small area of engagement minimizes the amount of friction and wear caused by the up and down movement of the flat. In this regard, the small engagement area also advantageously requires less oil to keep the surfaces properly lubricated.\nAs mentioned, the invention contemplates a predetermined flow path for directing lubrication to the roller member, and specifically, the bearing assembly thereof. The predetermined flow path, which in the preferred and illustrated form includes internal oiling channels formed in the valve lifter body that extend between the oil receiving area on the lifter body and the roller member, avoids the need to add oil squirters or add direct feed oil lines. This is desirable because oil squirters are not practical for use in aluminum and high performance steel connecting rods due to the loss of strength and stress riser resulting from the addition of the hole. Furthermore, the amount of the oil thrown from the squirters decrease as engine speeds decrease and are thereby inefficient if not unreliable.\nAlternatively, an external oiling channel can be provided on the surface of the lifter body. This external oiling channel is used to direct oil received by the oil receiving area, which is more preferably, an annular, circumferential groove about the lifter body that intersects the common oil passageway as the lifter reciprocates in the bore. As oil is received in the groove, the external oiling channel directs oil towards the housing portion of the lifter body where it may lubricate the roller and bearing assembly situated therein. Another advantage of using an annular groove and external oiling channel is that any oil thrown on the body of lifter may also be contained by the groove and channel and directed to the roller and bearing assembly.\nThe oil receiving area is in one form a transverse, through passageway that can be modified by adding of at least two round or oval shaped receiving areas on each side of the lifter body. The oil receiving areas are oriented perpendicular to the rolling direction of the roller, are ramped into the body and intersect the common transverse oil passageway in the engine block which feeds oil to the lifter galley. In the lifter body these oil receiving areas or inlets are connected by a passageway which travels through the body and parallel with the shaft of the roller. Also, additional inlets may be added to the front and back surfaces of the lifter body and connected to the internal passageway to feed additional oil into the passageway. Internal oiling channels have been added in the lifter body to direct the oil feed into the inlets and passageway. The oiling channels originate at the passageway and axially direct the oil through the body to the housing mounting the roller. To increase bearing and shaft life, at least two oiling channels are positioned to deposit oil between the housing and the outward sides of the roller to facilitate lubrication the shaft and bearing and to indirectly the surface of the roller. Additional oiling channels may be added to directly feed oil to the surface of the roller to directly lubricate the roller and camshaft lobes.\nTo prevent rotation of the valve lifter as it rapidly reciprocates up and down, a small guide or anti-rotation member has been added and fixed to the engine block in the lifter galley. The anti-rotation guide can span across two adjacent lifters and has a substantially flat main portion that sits on top of the engine block outside the bores. A tab extends perpendicular from the middle of the guide and in one form has a slot where a fastener may be inserted and threaded into the block of the engine to hold the guide stationary. Other methods of securing the anti-rotation guide to the block may also be employed. Each end of the guide that spans a lifter bore contains a small, crescent-shaped portion which depends from the main portion to form a shoulder therewith. The small crescent portion extends into the lifter bore with the curved portion of the crescent-shaped portion matching the curvature of the lifter bore to provide secured and flush engagement between the bore walls and the top of the block and the crescent-shaped portion of the anti-rotation member. The portion also has a planar bearing surface that mates with the front surface of the lifter for preventing rotation of the body of the lifter in the bore. To increase stability and decrease friction and wear on the valve lifter as it reciprocates, the lifter body has been modified by machining a short, planar surface on the front of the lifter. Due to the small contact surface created by the crescent-shaped portion of the anti-rotation guide and only a small portion of the lifter body need be planar. Now, as the lifter reciprocates in the bore, the front planar surface slides across the small planar surface of the anti-rotation guide containing its movement.\nThus, this guide provides an alternative to link bars which not only add excess material to the lifter assembly, but also present the potential for damage to the engine as the bars and attachment members wear due to the constant motion of the assembly. Furthermore, the small contact area created by the crescent-shaped portion minimizes friction and heat created thereof. Also, the guide also allows the mechanic to remove a single valve lifter from the engine by loosening the fastener and lifting and sliding the guide to allow the lifter to clear the guide; conventional link bars require the removal of the lifters as a pair. The capability to remove one lifter at a time is advantageous in engines where the pushrods may be of different lengths for the exhaust and intake valves. The mechanic needs to remove only one valve lifter and pushrod and thereby prevents the inadvertent switching of the pushrods during reassembly."}
{"text": "Since expanded resin beads are excellent in lightness in weight, cushioning properties, and heat resistance and high in a degree of freedom of shape design with which three-dimensional molded articles are obtained by means of in-mold molding, they are utilized as a cushioning material, a container, a heat insulating material, a vibration-damping material, or the like in multipurpose fields inclusive of packaging fields, commodities for living, building and civil engineering materials, vehicle members, and the like. The expanded resin beads are roughly classified into those made of a styrene-based resin as a base material resin and those made of an olefin-based resin as a base material resin. Among them, the expanded resin beads made of a styrene-based resin as a base material resin are used more frequently than the expanded resin beads made of an olefin-based resin as a base material resin for the reasons that the former is excellent in a balance between lightness in weight and compression strength, is easy for in-mold molding, and is inexpensive. But, though the expanded beads made of a styrene-based resin have an excellent aspect as described above, they are insufficient in heat resistance, chemical resistance, durability, toughness, and creep resistance characteristic depending upon uses.\nFor that reason, among the olefin-based resins, the expanded beads made of a propylene-based resin, which are excellent in heat resistance, chemical resistance, durability, toughness, creep resistance characteristic, and so on, are watched. But, with respect to in-mold molding using expanded propylene resin beads, there is involved difficulty in the aspect of molding processing, for example, a molding pressure of a heating medium, such as steam, etc., is high due to crystallinity or heat resistance of the propylene-based resin, and hence, an improvement is required. In addition, from the standpoint of performances of an expanded beads molded article, the requirements for lightness in weight and high impact energy absorbing performance are increasing, and a response to a more enhancement of rigidity of molded article may be considered.\nAs the conventional art, for example, PTL 1 describes that in order to utilize fusion characteristics by a propylene-based resin having a low melting point, thereby decreasing a molding pressure of steam at the time of in-mold molding, while appearance, heat resistance, and mechanical physical properties of a propylene-based resin having a high melting point, the propylene-based resin having a high melting point and the propylene-based resin having a low melting point are mixed under specified conditions. In addition, for example, PTLs 2 and 3 disclose that in order to decrease a molding pressure of steam at the time of in-mold molding of expanded propylene resin beads, expanded propylene resin beads in which a core layer thereof is covered with a resin having a low melting point are used."}
{"text": "In a copending U.S. patent application Ser. No. 526,700 entitled, \"Progressive Scan Television Display System Employing Interpolation in the Luminance Channel\" filed concurrently herewith in the name of D. H. Pritchard, a progressive scan television system is described. According to the Pritchard application, in the case of NTSC, 525 lines of the display are displayed in 1/60th of a second wherein alternate \"real\"0 and \"interpolated\" lines are successively displayed at some multiple (i.e., two-times) standard horizontal rate. During the next 1/60th of a second another 525 lines are displayed to complete a total frame time in 1/30th of a second, however, these second 525 lines are related to the first 525 lines such that in successive fields \"interpolated\" and \"real\" lines are displayed on top of each other. This progressive scan format results in the elimination of artifacts of \"interline flicker\" and \"line break-up with motion\" that exists in conventional two-to-one interlaced displays such as the NTSC system. The subjective effect of progressive scan is a flicker free, \" smooth\" or \"quiet\", picture presentation that is more pleasing to the viewer. The Pritchard application applies as well to other interlaced systems such as the 525/25 PAL system.\nTo perform progressive scanning as described in the aforementioned Pritchard application a speed-up processor is used for providing two horizontal lines of video during the time period of one horizontal line of incoming broadcast video signal. This speed-up processing requires that a line of video be written into memory in real time at a rate sufficient to sample the incoming video, i.e., at more than the Nyquist rate, and be readout of memory at some multiple of the input sample rate (generally, two-times). Such processing requires extremely high speed memories which are capable of operating at high sampling speeds, for example, 28 MHz for NTSC signals sampled at four-times color subcarrier. It is desirable to provide inexpensive speed-up processors which are capable of operating at these speeds without the use of memory elements which operate at extremely high speeds."}
{"text": "The methods described in this section could be practiced, but have not necessarily been previously conceived or pursued. Therefore, unless otherwise indicated herein, the methods described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nComputer networks that use routers, switches and other network elements are commonly managed using network management software systems. Examples of network management systems include Resource Management Essentials from Cisco Systems, Inc., San Jose, Calif., HP OpenView from Hewlett-Packard Company, Palo Alto, Calif., and others. Such network management systems are commonly used to support detection of network scenarios, including device faults, and to apply corrective instructions or configuration commands to network elements to address the scenarios.\nOne of the shortcomings of these network management software systems is the lack of guidance and automation for users of the systems. As a result, the network operator is required to know how to identify events that may represent symptoms of problems, how to determine whether symptoms actually represent the problems, how to diagnose the problems, how to select corrective action, and how to apply or perform the corrective action. In addition, the operator is also required to know the syntax of the commands and the usage of the commands in precise sequences to perform the aforementioned steps. All this knowledge that the operator mentally retains is neither easily transferable to another nor easily updateable to reflect platform or software changes.\nAnother shortcoming of these network management software systems is the primitive and proprietary interfaces of the systems. Specifically, many such systems support command line interfaces that accept character-based commands that conform to a complex grammar. In some cases, diagnosing faults requires the network operator to manually enter one command at a time. Numerous commands may be needed to determine a particular fault.\nAlso, each of these software systems is designed to handle a specific problem or scenario and has a distinct interface and feature set. As a result, the network operator is forced to learn multiple interfaces, feature sets, and subsequent updates or modifications to effectively manage multiple network scenarios.\nBased on the foregoing, there is a clear need for improved network management methods or systems that overcome the stated drawbacks of the current systems. There is a need to flexibly and effectively manage multiple types of devices and respond to various network scenarios. There is also a need to aggregate and interpret relevant information from different network devices."}
{"text": "RF (radio frequency) power architectures within the telecommunications field focus on achieving high DC-to-RF efficiency at significant power back off from Psat (the average output power when the amplifier is driven deep into saturation). This is due to the high peak to average ratio (PAR) of the transmitted digital signals such as W-CDMA (wideband code division multiple access), LTE (long term evolution) and WiMAX (worldwide interoperability for microwave access). The most popular power amplifier architecture currently employed is the Doherty amplifier. The Doherty amplifier employs a class AB main amplifier and a class C peaking amplifier, and efficiency is enhanced through load modulation of the main amplifier from the peaking amplifier. However, if high efficiency at a high output backoff (OBO) is required, a highly asymmetric ratio between the main and peaking amplifiers is required.\nThe Doherty architecture has an inherent degradation in the efficiency between the peak OBO point and the peak power point. To overcome this, a three way Doherty architecture can be used, in which the main class AB amplifier is replaced with a Doherty amplifier and load modulation is provided to the first peaking amplifier between the peak OBO point and the peak power point. However, the main amplifier is connected to the external load impedance (typically 50 Ohms) through a series of three ¼λ (quarter wavelength) transmission lines prior to any device impedance matching. This can lead to the amplifier being narrow band in nature due to the band-limiting characteristics of the ¼λ transmission lines. As such, three way Doherty amplifiers are typically designed for a specific band of operation used for wireless communication applications like WCDMA, LTE, WiMAX, etc. Such bands of operation are 1805-1880 MHz, 1930-1990 MHz, etc."}
{"text": "Chemically fexofenadine hydrochloride is 4-[4-[4-(hydroxydiphenylmethyl)-1-piperidinyl]-1-hydroxybutyl]-α,α-dimethylbenzene acetic acid hydrochloride. It is also known as terfenadine carboxylic acid metabolite. It is represented by Formula 1.\n\nFexofenadine hydrochloride is useful as an antihistamine, and does not cause the adverse effects associated with the administration of terfenadine including abnormal heart rhythms in some patients with liver disease or patients who also take the antifungal drug ketoconazole or the antibiotic erythromycin.\nU.S. Pat. No. 4,254,129 (“the '129 patent”) entitled Piperidine Derivatives issued on Mar. 3, 1981. The '129 patent relates to substituted piperidine derivatives and methods of making and using them. The disclosed compounds, including fexofenadine and its pharmaceutically acceptable salts and individual optical isomers, are purported to be useful as antihistamines, antiallergy agents and bronchodilators.\nThe '129 patent discloses a process for the preparation of fexofenadine having a melting point of 195-197° C. The recrystallization process exemplified therein in Example 3, column 13, involves use of a mixture of solvents for preparation of fexofenadine.\nWO 95/31437 discloses processes for preparing hydrated and anhydrous forms of piperidine derivatives, polymorphs and pseudomorphs thereof, which are useful as antihistamines, antiallergic agents and bronchodilators.\nWO 95/31437 discloses the preparation of anhydrous forms of fexofenadine hydrochloride by subjecting the hydrated fexofenadine hydrochloride to an azeotropic distillation or to water minimizing recrystallization. In the invention described in this application, unlike the process described in WO 95/31437, hydrated Fexofenadine Hydrochloride is not converted to anhydrous Fexofenadine Hydrochloride, but instead Fexofenadine is converted to Form A of Fexofenadine and then to anhydrous Form X of Fexofenadine Hydrochloride. The novel anhydrous crystalline form of Fexofenadine Hydrochloride is obtained according to the present invention directly from the novel precursor i.e. Fexofenadine without generating a hydrated form. The starting material used Fexofenadine (Base) is different than described in WO 95/31437.\nWO 00/71124A1 discloses amorphous fexofenadine hydrochloride process, its preparation and a composition containing it.\nFexofenadine obtained in the prior art processes, is a mixture of regioisomers of fexofenadine containing 33% of para isomer and 67% of meta isomer. These components are referred to as inseparable and it is also stated that it is not possible to obtain either of the regioisomers in substantially pure form. On the other hand, Fexofenadine prepared according to the process of this invention has a purity of >99.5%. In the novel crystalline Fexofenadine of this invention, the meta isomer of Fexofenadine is at a level of below 0.1%. Purity of fexofenadine is critical when it is used for the conversion to its hydrochloride salt since it is very difficult to remove any undesired impurities, including regioisomers, from the desired compound in last late processing stage. Removing the impurities increases the cost of production. Hence it is generally preferred that the HPLC purity of fexofenadine is greater than 99.5%.\nAnother beneficial aspect of the present invention is that, the fexofenadine hydrochloride is obtained in almost quantitative yield from the precursor i.e. fexofenadine. Almost quantitative yield means that the pure fexofenadine is converted to fexofenadine hydrochloride quantitatively (>92% yield of theory), with almost no yield loss, as the fexofenadine base itself is >99.5% pure as compared to fexofenadine prepared by the prior art processes."}
{"text": "The present invention relates to an image outputting system, which includes an image outputting apparatus and a controlling apparatus, which is coupled to the image outputting apparatus so as to develop first-type image data into second-type image data serving as an output possible format in the image outputting apparatus and sends the second-format image data to the image outputting apparatus.\nIn recent years, the image outputting apparatus, such as, for instance, a color laser printer employing an electro-photographic method, a compound apparatus having functions of a color laser printer, a scanner, a copier, a facsimile, etc. (hereinafter, referred to as a color laser printer as a whole), has been improved in its capability of producing a high quality image, and further, in its capability of speedily outputting a high quality color image without increasing its cost. Further, various kinds of finishers to be coupled to the color laser printer have been devised, and as a result, it becomes possible for the color laser printer to create and output printed products in such a output manner that various kinds of processing, such as a stapling, etc., are applied to the printed products. Reflecting such the recent trends, the color laser printer has been employed for the outputting use of final printed products, for instance, in a small-lot printing field, etc.\nOn the other hand, according to the proliferation of the DTP (Desk Top Publishing), etc., the image data are created on a client terminal, serving as an external terminal device including a personal computer, etc., by conducting an editing operation to be executed on the imaging software. Such the image data are to be elements of the image forming operation (hereinafter, referred to as elementary data) and includes vector data (the first-type image data) and raster data (the second-type image data). The vector data are called vector graphics, which represent each of objects as the image shown in FIG. 11 by using an aggregation of depicting information including coordinates of points and parameters for equations of lines and surfaces coupling objects to each other, painting colors (in FIG. 2, represented by the gradation values of dot percent for each of YMCK colors), special effects, font information representing a shape of character and its size, etc. While, the raster data are called bitmap graphics, in which one raster line is constituted by a plurality of pixels arrayed at predetermined intervals (resolution) in a horizontal direction, and further, a plurality of raster lines are also arranged at predetermined intervals in a vertical direction, and each of pixels is represented by using depicting information represented by gradation values. Further, the setting information including the first predetermined conditions, such as, for instance, image forming conditions with respect to the color adjustment, the layout such as allotment of the images, etc., the second predetermined conditions, such as, for instance, outputting conditions with respect to the output mode, such as a proofing, a stapling, a folding, etc., and the number of output copies, and a name, a name of terminal device, date-and-time information as an attribute information, are attached to the elementary data.\nThen, the outputting operation of the printing products is conducted on the basis of the elementary data. Concretely speaking, at first, the controlling apparatus coupled to the color laser printer receives the elementary data and the setting information from the terminal device through the network, and then, applies the color adjustment processing and the allotment processing to the elementary data so as to develop (convert) them into the printing data serving as image data having an output possible format for the color laser printer (the second-type image data), namely, for instance, halftone dot image data in which dots are formed within a predetermined area as shown in FIG. 13 (pixels in 5×5 area in the drawing) so as to express light and shade. according to the area of the dots. In other words, the elementary data are converted to the halftone dot image data, based on the first predetermined conditions indicated in the setting information. Then, the color laser printer conducts the image forming operation of the halftone dot image data based on the second predetermined conditions to output the printing products.\nIncidentally, when conducting the abovementioned image outputting operation, to illuminate the waste of printing products due to a certain defect generated in the printed contents, such as, for instance, a color defect, etc., which is found after print outputting operations of plural copies are completed, there has been a function of “proofing” in which only one copy to be used for confirmation of the printed contents is printed. According to the proofing, it becomes possible to reduce the waste of printing products, since only one copy is outputted from the color laser printer in order to confirm the contents of the printed copy and then a plurality of copies are outputted from the color laser printer.\nAs a result of the proofing, however, sometimes, it becomes necessary to change the image forming conditions with respect to the color adjustment, the layout such as allotment of the images, etc., serving as the first predetermined conditions in regard to the developing operation. It is impossible, however, for the operator to directly modify the halftone dot image data in order to cope with the abovementioned changing operation. Accordingly, to cope with such the changing operation, the operator had to stop the outputting operation based on the halftone dot image data concerned, and had to move from the place where the color laser printer is installed to the place where the terminal device is installed, in order to conduct the operation for changing the image forming conditions, the operation for adding the changed image forming conditions to the elementary data, and the operation for sending the elementary data added with the changed image forming conditions to the controlling apparatus, at the terminal device, and then, had to again conduct the developing operation mentioned in the above at the controlling apparatus. As described in the above, according to the conventional method, it had taken much time and labor to change the image forming conditions.\nFurther, as a function of the proofing, for instance, Patent Document 1 sets forth the feature for conducting proofing operations with respect to all of the printing modes so as to make it possible to select a preferable printed result from them. However, this feature requires the developing operations and outputting operations in regard to all of the printing modes, resulting in a large amount of waste of time and expendable supplies.\n[Patent Document 1]    Tokkaihei 11-134147 (Japanese Non-Examined Patent Publication)"}
{"text": "1. Field of the Invention\nThe present invention relates to a lamp holder, a backlight device using the same, and a display using the same.\n2. Description of the Related Art\nFor example, a liquid crystal panel used for a liquid crystal display such as a liquid crystal television additionally requires a backlight device as an external lamp because it does not spontaneously emit light. The backlight device is placed on a backside of the liquid crystal panel, and broadly includes a base made of a metal with an open surface at the side of the liquid crystal panel, a number of cold cathode tubes housed in the base as lamps, and a number of optical members (diffusion sheets and the like) which are arranged in an open portion of the base to efficiently irradiate light, which is emitted by the cold cathode tubes, to the liquid crystal panel side, and includes lamp clips for holding the cold cathode tubes, each having a slim tubular shape with respect to the base.\nAn example of the lamp clips as described above can be found in Japanese Patent Laid Open No. 2001-210126. The lamp clips of this invention are made of a synthetic resin, and the invention includes a mounting plate which is applied to an inner surface of the base, a locking part which is protruded to the base side from the mounting plate, is inserted into a mounting hole of the base to be capable of being locked to its peripheral edge, and a lamp holding part which is protruded to the side opposite from the locking part from the mounting plate to be capable of holding the cold cathode tube so as to surround the peripheral surface of the cold cathode tube. The lamp holding part has a C-shaped sectional configuration which is opened upward, and is elastically deformable so as to open outward during attachment and detachment of the cold cathode tube.\nIncidentally, typically there is a variation in the thickness of the cold cathode tube that cannot be prevented due to a manufacturing error. Thus, setting the size of the lamp holding part which holds such a cold cathode tube can be problematic.\nIf the size of the lamp holding part is set with the thickest cold cathode tube in tolerance as a reference, the clearance which occurs between the lamp holding part and the cold cathode tube becomes too large when a relatively thin cold cathode tube is mounted, and large backlash occurs to the cold cathode tube.\nIf the size of the lamp holding part is set with the thinnest cold cathode tube in tolerance as the reference in contrast with the above-described situation, the force which is required for opening and deforming the lamp holding part becomes too large when a relatively thick cold cathode tube is mounted. Even if it is mounted, the elastic rebound force of the lamp holding part becomes too large, and therefore, there arises the problem that the cold cathode tube is urged in the detaching direction by the elastic rebound force and the cold cathode tube is easily detached, and in short, this setting method is not favorable with respect to attaching and detaching operability and holding performance. Nevertheless, there are limitations in the molding technique and strength in molding the lamp holding part to be thin in order to reduce the elastic rebound force of the lamp holding part, thus bringing about difficulties in coping with this situation."}
{"text": "With the increased popularity of mobile computing devices, there is an ever growing need for high capacity data throughput to wireless computing devices, such as hand held “smart” cellular telephones, wireless network adapters, and other types of wireless data devices. To meet this need, macrocell and picocell base stations are being widely installed. A picocell base station is typically a relatively low cost, small, reasonably simple unit that connects to a base station controller. A macrocell base station serves a larger geographic area than a picocell base station, but is much more expensive to install.\nBecause fiber optic cable IS not always available, or may be cost prohibitive, telecommunications carriers are beginning to rely more greatly on wireless infrastructure to deploy backhaul data capacity to macrocell and picocell base stations. One challenge when deploying wireless backhaul is not having a clear line of sight from a macrocell or picocell base station to a backhaul end-location. Structures such as buildings, roads, vegetation and residential homes may prevent point-to-point wireless connections. Moreover, weather and other factors might greatly impact current wireless backhaul solutions.\nIt is with respect to these and other considerations that the disclosure made herein is presented."}
{"text": "1. Field of the Invention\nThis invention relates to the field of ceramics and particularly to ZrO.sub.2 ceramics.\n2. Description of the Prior Art\nDuring cooling, ZrO.sub.2 undergoes a martensitic-type transformation from a tetragonal crystal structure to a monoclinic crystal structure with a concurrent increase in volume and an anisotropic shape change. For pure ZrO.sub.2 the transformation begins at about 1200.degree. C. and proceeds until complete at about 600.degree. C.\nAttempts have been made to utilize this transformation in order to improve the fracture toughness of ceramic composites. In one approach, ZrO.sub.2 particles have been added to an Al.sub.2 O.sub.3 matrix to form a second phase dispersion (N. Claussen, J. Am. Ceram. Soc. 59, pg. 49, 1976). Expansion and shape change of the ZrO.sub.2 as it transformed from the high temperature tetragonal structure to the room temperature monoclinic structure created microcracks. The resulting increase in fracture toughness was attributed to energy absorption by these microcracks.\nMore recently, attempts have been made to increase the toughness of ZrO.sub.2 ceramics by taking advantage of metastable grains of tetragonal ZrO.sub.2 within a surrounding matrix. These are grains of ZrO.sub.2 which are tetragonal rather than monoclinic despite the fact that their temperature is below the unconstrained equilibrium transformation temperature range.\nThe metastable condition can be obtained by surrounding fine grains of ZrO.sub.2 in a constraining matrix such as Al.sub.2 O.sub.3. The matrix constrains the volume and shape change associated with the transformation of the grains of ZrO.sub.2 and holds the ZrO.sub.2 in its tegragonal state.\nThe tetragonal grains of ZrO.sub.2 increase the fracture toughness of the ceramic composite by increasing the energy required for a crack to propogate. If a crack starts in the ceramic composite, the metastable grains of tetragonal ZrO.sub.2 in the stress field adjacent the crack transform to the stable monoclinic structure. The work done by the applied stresses to reduce this transformation is loss and thus the stress-induced transformation increases the material's fracture toughness.\nMetastable tetragonal grains of ZrO.sub.2 have been observed in an Al.sub.2 O.sub.3 /ZrO.sub.2 ceramic composite containing 17 volume % ZrO.sub.2 (N. Claussen, J. Am. Ceram. Soc. 59, pg. 85, 1978). However, to maintain the metastable tetragonal structure, the ZrO.sub.2 grains had to be less than about 0.5 .mu.m in diameter. Larger grains transformed to the stable monoclinic structure. Additional work has shown that the amount of metastable tetragonal ZrO.sub.2 that can be retained in the matrix decreases as the volume % of ZrO.sub.2 in the Al.sub.2 O.sub.3 /ZrO.sub.2 ceramic composite increases. Very little of the ZrO.sub.2 can be retained in the metastable tetragonal structure in Al.sub.2 O.sub.3 /ZrO.sub.2 composites having more than 20 volume % ZrO.sub.2. Such limitations of grain size and volume % of ZrO.sub.2 reduces the practicality and the toughness of prior art Al.sub.2 O.sub.3 /ZrO.sub.2 ceramic composites."}
{"text": "The present disclosure relates to a backlight assembly for a liquid crystal display, and specifically, to a backlight assembly with high efficiency and brightness."}
{"text": "In computer systems in general and in small systems in particular, it is desirable to reduce the total amount of hardware required for the system consistent with some specified level of machine performance. Data lines connecting various external devices such as keyboard terminals, CRT displays, credit card readers, etc., to a central processor would normally require separate I/O pin connections for each data line. In modern computers, large scale integrated circuitry (LSI) is widely used and, as is well known, these circuit modules are limited in the number of external circuit connections or pins which can be placed on an LSI module. Hence, by utilizing the same data lines for both input and output, the number of pins required is cut in half. In large, high speed computing systems, it generally has not been feasible to utilize the same data lines for both reading data into and out of the system, as many operations are held up pending the termination of an existing operation. While such waiting is necessary with a central memory or other shared functional unit, it is not so critical within the CPU as modern computers have extremely high speed memory circuits, etc., capable of performing required operations in extremely short periods of time. However, with smaller lower speed computers where both cost and manufacturing feasibility are primary factors, and time is somewhat secondary, it becomes possible to consider bidirectional data lines or half-duplex type of operations.\nWith the advent of modern integrated circuits and micro electronic technology, it has been possible to construct relatively small inexpensive computers utilizing separate source and destination buses within the computers which allows for considerably higher speed operations without greatly increasing the cost of construction. However, in the past, such two bus computers required large numbers of terminal pins for connecting the computer to both the input and output data lines. As stated previously, this causes severe pin connection problems when LSI is used in the host computer. It is the purpose of the present invention to provide a solution to the terminal or pin connection problem.\nIt has been found that satisfactory overall system performance is possible with a bidirectional data bus connected to a small two separate internal bus computer system by providing internal gating and control means for selectively connecting the I/O bus to the source and destination bus of said computer utilizing internal control circuitry together with a special adaptor connecting I/O devices to the I/O bus and for effecting requisite I/O control.\nIt has been further found that the requisite control functions can be efficiently and inexpensively performed by the use of micro programs stored in a read only memory having special decoding and control circuitry for effecting the desired gating operations."}
{"text": "In systems for controlling electrical processes, such as for controlling generation and transmission of electrical power, i.e. high-voltage power transmission and generation systems, it is necessary to measure the voltage at various points in the system, such as for instance in a power line in order to provide control and protection of the system.\nIn order to provide a reliable system there may be redundancy in relation to the control and protection devices provided. Such a redundancy means that several different parallel protection and control devices may need to receive measurements from the same element in the process.\nWO2006/128397 discloses an HVDC system where two separate connections to a power line are provided in order to provide redundancy in relation to measurement results.\nHowever, when the system is a high-voltage system it is often necessary to limit the number of contacts to the system element to as few as possible. There may otherwise be a risk for flashover. This is especially the case in Ultra high HVDC (High Voltage Direct Current) power transmission systems.\nOne way to solve this problem may be to provide duplicated amplifiers connected to a voltage dividing element. Voltage measurements of a system element, such as a power line may be performed through duplicated amplifiers directly connected to such a primary voltage dividing element. However, with this solution it may not be possible to ensure a high enough reliability. The amplifiers may furthermore not be possible to be replaced during operation.\nThere is therefore a need for an improvement when measuring voltages of electrically conducting elements in a system for controlling an electrical process."}
{"text": "With the advent of GPS technology, it has become possible to determine the precise location of any GPS receiver on the globe. GPS technology is often combined with mapping software in order to visually indicate the location of a GPS receiver on a map. Some of the portable devices that are currently available on the market include a screen for displaying a map along with a visual indicator that represents the GPS location of the device, however, there are several drawbacks associated with these devices.\nDownloading geographic maps onto the portable devices from a network is very time consuming and there are often licensing fees associated with map use. In addition, the size of the screen is typically minimized in order to reduce the overall size and weight of the portable device. This may result in a user spending valuable time attempting to understand the map if insufficient information is provided on a single screen.\nAs with any map, time is also spent becoming oriented with one's surroundings and correlating them to the map on the screen. If the user is not particularly map literate, it may take a very long time for the user to determine his or her location. Further, in a region that is poorly mapped or a region that includes few roads or landmarks, mapping one's GPS location on a portable device may not be very useful."}
{"text": "1. Field of the Invention\nThe present invention relates to a skin-cleansing composition in the form of an oil-in-water (O/W) emulsion, which may be used for cleansing the skin of the human face, neck and/or body. In particular, this composition, which enables dead cells to be removed gently from the skin, constitutes a \"3 in 1\" product which is particularly suitable for sensitive skins.\nThe invention also relates to a process for cleansing the skin of the body and/or face.\n2. Discussion of the Background\nExfoliant or scrubbing products, also referred to as \"scrubs\", contain exfoliant particles which consist of abrasive materials such as polyethylene powder, walnut shell powder or apricot or almond powder. Unfortunately, these scrubbing products are irritating due to the very presence of these particles which rub the skin without melting and which remain on its surface. On account of their irritant nature, it is not possible to use these scrubbing products daily; depending on the sensitivity, there must be longer or shorter time intervals between uses thereof, in particular of several days, or even of several weeks for very sensitive skin. Moreover, in order to overcome skin irritation, it is necessary, after use, to apply a care cream which provides protection and moisturization to the skin.\nThere thus is a need for a scrubbing-product which is comfortable to use, which does not contain abrasive exfoliant particles and which does not require the subsequent application of a care cream.\nThe composition according to the present invention overcomes the drawbacks described above. Applicant has found, surprisingly, that it is possible to have a \"3 in 1\" scrubbing product, that is to say a product which simultaneously cleanses, scrubs and cares for the skin, containing no abrasive particles, by providing an oil-in-water emulsion comprising lipid grains filled with a cleansing agent."}
{"text": "The invention relates generally to electronic capacitors, and more specifically to foil stacked capacitors using aluminum and polymer.\nElectrical circuits often include capacitors for various purposes such as filtering, bypassing, power decoupling, and to perform other functions. High-speed digital integrated circuits such as processors and computer chipsets in particular typically perform best when the power supplied to the integrated circuit is filtered with a capacitor placed physically close to the integrated circuit.\nSuch power decoupling capacitors function to smooth out irregularities in the voltage supplied to the integrated circuits, and so serve to provide the integrated circuits with a more ideal voltage supply.\nBy placing the decoupling capacitors near the integrated circuit, parasitic impedances such as printed circuit board path resistance or inductance are minimized, allowing easy and efficient transfer of energy from the.decoupling capacitor to the integrated circuit. Minimization of series resistance and inductance in the capacitor itself is also desirable for the same purposes, and results in a more efficient and desirable decoupling or bypass capacitor.\nThe internal series resistance of the capacitor is typically known as the Equivalent Series Resistance, or ESR. Similarly, internal series inductance is known as Equivalent Series Inductance, or ESL. Both of these parameters can be measured for a given capacitor, and are among the basic criteria used to select capacitors for applications such as integrated circuit power supply decoupling.\nPast efforts to minimize ESL and ESR have included solutions such as using multiple types of capacitors in parallel or combination series-parallel configurations, configured to product the desired capacitance at the very low ESR and ESL levels required. For example, tantalum capacitors in the order of 4.7 uF in parallel with 0.01 uF ceramic chip capacitors were often sufficient for lower-speed digital logic circuits of previous decades. But, new high speed digital logic circuits such as high-performance computer processors require both greater capacitance because of the amount of power dissipated, and lower ESR and ESL because of the very high speeds at which the processors operate.\nIt is also desirable for capacitors to have a physically small size, so that they do not take an unduly large amount of printed circuit board space. This is why space efficient capacitor technologies such as tantalum and electrolytic capacitors are often implemented in circuits despite typically having relatively high inductance, resistance, dielectric absorption, and other unfavorable characteristics. Mitigation of unfavorable capacitor characteristics of electrolytic or tantalum capacitors often also requires use of parallel capacitors with more favorable characteristics as secondary or supplemental decoupling capacitors.\nWhat is desired is a single capacitor design that provides low ESR and ESL with large capacitance, and that is physically compact."}
{"text": "Wireless networks are based on transmission and reception of Radio Frequency (RF) electro-magnetic waves between at least two devices. During operation, the RF connection will experience a wide range of attenuations due to e.g., atmospheric effects, intervening obstacles, etc. Generally, it is assumed that the radio link between two devices is substantially symmetric; i.e., transmit and receive RF links for a given device are substantially similar in performance and attenuation. In symmetric link operation, each device assumes that their perceived performance is representative of a peer device's perceived performance at the other end of an RF link.\nHowever, empirical evidence suggests that certain device usage scenarios can create asymmetric attenuation in each direction of an RF link, e.g. in transmit and receive RF links for a given device. For example, a user operating a wireless communication device (such as e.g., the iPad™ developed and manufactured by the Assignee hereof) in communication with a nearby Long Term Evolution (LTE) cellular network can experience a wide range of radio link performance based on e.g., the positioning of the wireless communication device with respect to the user's body. Specifically, the wireless communication device's RF transceiver may be configured to limit transmit power when the transmitter is next to the user's body (or adjacent to another object). The degree of transmit power reduction or attenuation at the wireless communication device can be based on, in one case, regulatory concerns, which can dictate Specific Absorption Rate (SAR) requirements permitted for emitted RF energy by the wireless communication device.\nIn the foregoing scenario, only the transmit link may be affected by the regulatory provisions, which can limit operating conditions and therefore influence performance; the receive link can remain fully operational. In some situations, the wireless communication device may not be able to successfully interact with a wireless network because its transmit power may be limited, and transmissions may be necessary to establish and/or to maintain an RF link with the wireless network, (i.e., the wireless communication device can be unable to set up the link by only receiving commands or data from a wireless network.) To make matters worse, since reception quality at the wireless communication device can be acceptable in this scenario, the wireless communication device may be unable to attempt to invoke remedial or corrective action; e.g., execute a handover to another LTE evolved NodeB (eNB), perform a cell selection/reselection, etc. Existing wireless communication devices may therefore be caught in a marginal and/or unusable operational state, thereby degrading performance and user experience.\nAccordingly, improved methods and apparatus are needed for handling radio link imbalances such as for example those described in the foregoing scenario."}
{"text": "1. Field of the Invention\nThe present invention relates to the prescription medical treatment for sinusitis, more specifically, to a medicinal package that improves compliance with the treatment regimen prescribed for sinusitis.\n2. The Prior Art\nSinusitis is a common disorder affecting an estimated 10% of the United States population and affecting all age groups, including children and the elderly. The problem is increasing in prevalence and in 1994, sinusitis accounted for 25 million office visits to physicians in the United States. Sinusitis can be defined as an inflammation of the paranasal sinuses which manifests as a purulent (infected) nasal discharge, nasal congestion, pain in the sinus areas (cheeks, forehead, around eyes, sides of nose), which may be associated with fever, headache, dental pain, earache, post-nasal discharge, cough, sore throat, conjunctival inflammation, foul breath, and olfactory loss. Its temporal manifestations vary from an acute illness of less than three weeks, to recurring episodes, to an unremitting chronic condition. Complications of inadequately treated sinusitis, in addition to chronicity, can be grave because of the proximity of the sinuses to the bony walls enclosing the eyes (orbits) and brain, and include orbital celulitis, optic neuritis, cavernous sinus thrombosis, epidural and subdural infection, meningitis, cerebritis, brain abscess, blindness, and even death.\nThe management of sinusitis is predicated upon what is known of the pathophysiology of this disorder. The paranasal sinuses consist of a series of bony pouches adjacent to the nasal cavity in the frontal, maxillary, ethmoid and sphenoid regions, which are lined by pseudostratified, ciliated epithelium. Mucous is produced by epithelial goblet cells and submucosal seromucous glands. The blanket of mucous covering the epithelial surface of the sinuses is moved in an orderly fashion by cilia towards natural ostea which lead into the nasal cavity, thereby allowing constant drainage of the sinuses. When the flow of mucous from the sinuses is interrupted, the retained secretions become thickened, the adjacent mucous membranes become inflamed and both mucous and sinus membranes are subject to infection.\nPharmacotherapy for sinusitis is therefore directed at:\n(1) reestablishing patency of the sinus ostea (openings),\n(2) reestablishing the orderly flow of mucous, and (3) treating the infection. These three objectives conventionally require multiple prescriptions of individual medications, with a typical regimen including: (1) an oral decongestant to shrink the swelling of the sinus membranes thereby opening the sinus exit pathway, (2) an expectorant to increase respiratory tract fluid secretions, reduce their viscosity, and increase the efficacy of the mucociliary mechanism and facilitate mucous flow, and (3) an antibiotic to treat the infection.\nThe choice of medications and their use together are dependent on numerous considerations besides the mechanism of action and risks of the individual medications. These considerations include absorption, time of onset after dosing, rate of elimination, duration of action after dosing, therapeutic effect by virtue of combination, and side effects by virtue of combination. Medication error and misuse due to a multiplicity of medications and modalities pose an additional risk. Medical and pharmaceutical expertise is clearly required to formulate a treatment regime utilizing a combination of medications and appropriate instructions for use by a lay individuals affected by sinusitis.\nSuccess of such a treatment regimen is contingent upon compliance for a 10-14 day period for acute sinusitis and a 3-8 week period in children and individuals with chronic sinusitis. Previous compliance studies have demonstrated three important considerations which adversely affect compliance: (1) increased complexity of the treatment regimen, (2) poor patient understanding of the treatment rationale, and (3) difficulty of use. Indeed, the multiplicity of medications necessary for sinusitis treatment increases the complexity of the regimen, patients may not fully understand the benefit of each component, and the convention of multiple containers and separate instructions for each component make complying with the regimen more difficult.\nUnited States health care experts conservatively estimate that half of the 1.8 billion prescription medications dispensed yearly are not taken as prescribed. Because of its potentially negative consequences, many consider lack of compliance with treatment regimens to be one of the most serious problems facing health care today. The multiplicity of medications necessary for effective sinusitis treatment makes it especially susceptible to non-compliance.\nSolutions to the compliance problem have been put forth by others. Typical of such solutions is the compartmented pillbox, where the medications are stored in compartments representing times of the day and different days. The major shortcoming of the compartmented pillbox is that the patient still receives the medications in separate containers and then must sort the various medications and store them in the proper compartments in the pillbox. This can be a complex and difficult task, especially when the medications are similar in appearance. And there is no guarantee that the medications will be sorted and stored correctly.\nAn object of the present invention is to provide a means for increasing compliance with medication regimens for treating sinusitis.\nAnother objective is to provide a sinusitis patient with a unified, understandable, and organized treatment regimen for sinusitis.\nA further object is to minimize complexity and facilitate ease of use of a sinusitis treatment regimen.\nThe preferred embodiments of the present invention comprise a multiplicity of medications for sinusitis physically arranged so as to simplify their use, functional indicia and instructions for coordinating the medications together as a regimen, and unification of these elements within a pharmaceutical dispensing assemblage.\nWith the present invention, all of the medications for the treatment regimen are prepackaged into a single prescription package for the patient. The patient only deals with a single package, rather than the multiplicity of packages of the prior treatment regimens. The medication is organized into event modules associated with daily events at which the medication is taken. The event may be a time of day or an activity that is performed during the day. Indicia representing the events associated with the event modules lead the patient clearly through the treatment regimen over its full time period, leading to a greater degree of compliance with the regimen and a greater probability that the treatment will be successfully completed.\nThe medication dosages are stored in either blister packs or pouches. The blister pack includes a clear plastic sheet with pockets for the dosages and a rupturable or pealable cover for retaining the dosages in their pockets until manually removed. The pouch is a bag composed of thin sheets of plastic or foil and is typically opened by tearing. The present invention can be used with many physical forms of medication, but the preferred forms are those that are most easily taken, such as tablets, capsules, and liquid-gels.\nThere are two basic preferred embodiments of the present invention, the box embodiment and the card embodiment. The box embodiment includes a box and a plurality of event modules. Each event module is either a blister pack or a pouch and is identified by an event indicia. A set of one day\"\"s worth of event modules may be physically combined into a day group. The event modules are organized within the box to present the treatment regimen in a logical progression. In one form, the box has dividers that define compartments, where all of the event modules for one event reside within one associated compartment. In another form, the box is tall, with a slot on one side at the bottom through which one event module fits. The event modules are stacked within the container in chronological order or are all connected together and rolled into a loop. Each event module is removed from the box when needed by sliding it out of the slot.\nThe card embodiment includes a number of medication dosages in a blister pack container organized into day modules and event modules. Each day module represents a single day of the treatment regimen and includes one of each type of event module. The day modules are arranged in single or multiple rows or columns. All of the event modules of a single row or column that are defined by the same event are arranged in a continuous line. Each day module includes a day indicia indicating the day of the treatment regimen that the dosages of that day module are to be taken and each event module line is associated with an event indicia. Optionally, the day modules are delimited by perforations that allow the manual separation of a day module from the card.\nIn all embodiments, the assemblage includes an instruction area which contains any information deemed necessary to the safe use of the medications. Such information includes, but is not limited to, a graphical depiction of each event module, a graphical medication legend, and instructions for use.\nOther objects of the present invention will become apparent in light of the following drawings and detailed description of the invention."}
{"text": "This invention relates in general to sample inspection systems and, in particular, to an improved inspection system with good sensitivity for particles as well as crystal-originated-particles (COPs). COPs are surface breaking defects in semiconductor wafers which have been classified as particles due to inability of conventional inspection systems to distinguish them from real particles.\nSystems for inspecting unpatterned wafers or bare wafers have been proposed. See for example, PCT patent application Ser. No. PCT/US96/15354, filed on Sep. 25, 1996, entitled “Improved System for Surface Inspection.” Systems such as those described in the above-referenced application are useful for many applications, including the inspection of bare or unpatterned semiconductor wafers. Nevertheless, it may be desirable to provide improved sample inspection tools which may be used for inspecting not only bare or unpatterned wafers but also rough films. Another issue which has great significance in wafer inspection is that of COPs. These are surface-breaking defects in the wafer. According to some opinions in the wafer inspection community, such defects can cause potential detriments to the performance of semiconductor chips made from wafers with such defects. It is, therefore, desirable to provide an improved sample inspection system capable of detecting COPs and distinguishing COPs from particles."}
{"text": "The present invention relates in general to implantable, cardiac electrical stimulation devices such as pacemakers or implantable cardioverter-defibrillators. In particular, this invention pertains to a system and method for detecting ventricular capture using far-field sensing of the ventricular R-wave.\nImplantable medical devices, such as pacemakers, defibrillators, cardioverters, and implantable cardioverter-defibrillators (xe2x80x9cICDsxe2x80x9d), collectively referred to herein as implantable cardiac stimulating devices, are designed to monitor and stimulate the heart of a patient who suffers from a cardiac arrhythmia. Using leads connected to a patient\"\"s heart, these devices typically stimulate the cardiac muscle (myocardium) by delivering electrical pulses in response to measured cardiac events that are indicative of a cardiac arrhythmia. Properly administered therapeutic electrical pulses often successfully reestablish or maintain the heart\"\"s regular rhythm.\nImplantable cardiac stimulating devices can treat a wide range of cardiac arrhythmias by using a series of adjustable parameters to alter the energy, shape, location, and frequency of the therapeutic pulses. The adjustable parameters are usually defined in a computer program stored in a memory of the implantable device. The program, which is responsible for the operation of the implantable device, can be defined or altered telemetrically by a medical practitioner using an external implantable device programmer.\nConventional programmable cardiac stimulation devices are generally of two types: (1) single-chamber, or (2) dual-chamber. In a single-chamber pacemaker, the pacemaker provides stimulation pulses to, and senses cardiac activity within, a single-chamber of the heart, either the right ventricle or the right atrium. In a dual-chamber pacemaker, the pacemaker provides stimulation pulses to, and senses cardiac activity within, two chambers of the heart, e.g., both the right atrium and the right ventricle. The left atrium and left ventricle can also be sensed and paced, provided that suitable electrical contacts are effected therewith. The recent development of multi-chamber cardiac stimulation devices allows sensing and pacing in up to all four chambers of the heart. Recent clinical evidence suggests multi-chamber stimulation may have important hemodynamic benefit in patients suffering from heart failure and may be effective in preventing arrhythmias in patients prone to sustained or frequent arrhythmias.\nCardiac stimulation devices have a great number of adjustable parameters that must be tailored to a particular patient\"\"s therapeutic needs. One adjustable parameter of particular importance is the output stimulation energy. For the stimulation device to perform its intended function, it is critically important that the delivered electrical stimuli be of sufficient energy to depolarize the cardiac tissue, a condition known as xe2x80x9ccapturexe2x80x9d.\nWhen a pacemaker stimulation pulse stimulates either the atrium or the ventricle during an appropriate portion of a cardiac cycle, it is desirable to have the heart properly respond to the stimulus provided. Every patient has a xe2x80x9ccapture thresholdxe2x80x9d which is generally defined as the minimum amount of stimulation energy necessary to effect capture. Capture should be achieved at the lowest possible energy setting yet provide enough of a safety margin so that, should a patient\"\"s threshold increase, the output of an implantable stimulation device, i.e. the stimulation energy, will still be sufficient to maintain capture. Dual-chamber and multi-chamber stimulation devices may have differing atrial and ventricular stimulation energy that correspond to the capture thresholds of the targeted cardiac chamber.\nThe earliest pacemakers had a predetermined and unchangeable stimulation energy, which proved to be problematic because the capture threshold is not a static value and may be affected by a variety of physiological and other factors. For example, certain cardiac medications may temporarily raise or lower the threshold from its normal value. In another example, fibrous tissue that forms around stimulation electrodes within several months after implantation may raise the capture threshold.\nAs a result, some patients eventually suffered from loss of capture as their pacemakers were unable to adjust the pre-set stimulation energy to match the changed capture thresholds. One attempted solution was to set the level of stimulation pulses fairly high so as to avoid loss of capture due to a change in the capture threshold. However, this approach resulted in some discomfort to patients who were forced to endure unnecessarily high levels of cardiac stimulation. Furthermore, such stimulation pulses consumed extra battery resources, thus shortening the useful life of the pacemaker.\nWhen programmable pacemakers were developed, the stimulation energy was implemented as an adjustable parameter that could be set or changed by a medical practitioner. Typically, such adjustments were effected by the medical practitioner using an external programmer capable of communication with an implanted pacemaker via a magnet applied to a patient\"\"s chest or via telemetry. The particular setting for the pacemaker\"\"s stimulation energy was usually derived from the results of extensive physiological tests performed by the medical practitioner to determine the patient\"\"s capture threshold, from the patient\"\"s medical history, and from a listing of the patient\"\"s medications. While the adjustable stimulation energy feature proved to be superior to the previously known fixed energy, some significant problems remained unsolved. In particular, when a patient\"\"s capture threshold changed, the patient was forced to visit the medical practitioner to adjust the stimulation energy accordingly.\nTo address this pressing problem, pacemaker manufacturers have developed advanced stimulation devices that are capable of determining a patient\"\"s capture threshold and automatically adjusting the stimulation pulses to a level just above that which is needed to maintain capture. This automatic capture feature improves the patient\"\"s comfort, reduces the necessity of unscheduled visits to the medical practitioner, and increases the pacemaker\"\"s battery life by conserving the energy used to generate stimulation pulses.\nHowever, many of these advanced pacemakers require additional circuitry and/or special sensors that must be dedicated to capture verification. This requirement increases the complexity of the pacemaker system and reduces the precious space available within a pacemaker\"\"s casing, and also increases the pacemaker\"\"s cost. As a result, pacemaker manufacturers have attempted to develop automatic capture verification techniques that may be implemented in a typical programmable pacemaker without requiring additional circuitry or special dedicated sensors.\nA common technique used to determine whether capture has been effected is monitoring the patient\"\"s cardiac activity and searching for the presence of an xe2x80x9cevoked responsexe2x80x9d following a stimulation pulse. The evoked response is the response of the heart to the application of a stimulation pulse. The patient\"\"s heart activity is typically monitored by the stimulation device by keeping track of the stimulation pulses delivered to the heart and examining, through the leads connected to the heart, electrical signals that are manifest concurrent with depolarization or contraction of muscle tissue (myocardial tissue) of the heart. The contraction of atrial muscle tissue is evidenced by generation of a P-wave, while the contraction of ventricular muscle tissue is evidenced by generation of an R-wave (sometimes referred to as the xe2x80x9cQRSxe2x80x9d complex).\nWhen capture occurs, the evoked response is an intracardiac P-wave or R-wave that indicates contraction of the respective cardiac tissue in response to the applied stimulation pulse. For example, using such an evoked response technique, if a stimulation pulse is applied to the ventricle, a response sensed by ventricular sensing circuits of the stimulation device immediately following the application of the stimulation pulse is presumed to be an evoked response that evidences capture of the ventricle.\nHowever, it is for several reasons very difficult to detect a true evoked response. First, because the ventricular evoked response is a relatively small signal, it may be obscured by a high-energy stimulation pulse and therefore difficult to detect and identify. Second, the signal sensed by the stimulation device\"\"s sensing circuitry immediately following the application of a stimulation pulse may be not an evoked response but noise, such as electrical noise caused, for example, by electromagnetic interference, or myocardial noise caused by random myocardial or other muscle contraction.\nAnother signal that interferes with the detection of an evoked response, and potentially the most difficult for which to compensate because it is usually present in varying degrees, is lead polarization. A lead/tissue interface is that point at which an electrode of the stimulation lead contacts the cardiac tissue. Lead polarization is commonly caused by electrochemical reactions that occur at the lead/tissue interface due to application of an electrical stimulation pulse, such as a stimulation pulse, across the interface.\nBecause the evoked response is sensed through the same lead electrodes through which the stimulation pulses are delivered, the resulting polarization signal, also referred to as an xe2x80x9cafterpotentialxe2x80x9d, formed at the electrode can corrupt the evoked response that is sensed by the sensing circuits. This undesirable situation occurs often because the polarization signal can be three or more orders of magnitude greater than the evoked response. Furthermore, the lead polarization signal is not easily characterized; it is a complex function of the lead materials, lead geometry, tissue impedance, stimulation energy and other variables, many of which are continually changing over time.\nIn each of the above cases, the result may be a false positive detection of an evoked response. Such an error leads to a false capture indication, which in turn, leads to missed heartbeats, a highly undesirable and potentially life-threatening situation. In dual chamber and multichamber stimulation, successful capture all stimulated chambers is critical to maintaining proper synchrony of heart chamber contractions. Loss of optimal atrial-ventricular synchrony or inter-ventricular synchrony may have deleterious hemodynamic effects.\nAnother problem results from a failure by the stimulation device to detect an evoked response that has actually occurred. In that case, a loss of capture is indicated when capture is in fact present, also an undesirable situation that will cause the device to unnecessarily deliver a high-energy back-up stimulation pulse and invoke the threshold search function in a chamber of the heart.\nAutomatic threshold testing is invoked by the stimulation device when loss of ventricular capture is detected or on a predetermined periodic basis. An exemplary threshold test is performed as follows. When loss of capture is detected, the device increases the stimulation pulse energy to a relatively high predetermined testing level at which capture is certain to occur, and thereafter decrements the output energy until capture is lost. The stimulation energy is then set to a level slightly above the lowest output energy at which capture was still detected. Thus, capture verification is of utmost importance in proper determination of the stimulation energy.\nWhen a ventricular stimulation pulse is properly captured in the ventricle, a subsequent ventricular contraction results in an R-wave which may be sensed through an atrial lead, in patients with intact atrioventricular (xe2x80x9cAVxe2x80x9d) conduction, as a xe2x80x9cfar-fieldxe2x80x9d signal, also referred to herein as xe2x80x9cfar-field R-wavexe2x80x9d or xe2x80x9cfar-field evoked responsexe2x80x9d. The far-field R-wave confirms successful ventricular capture because the ventricular contraction only occurs after a properly captured ventricular stimulation pulse.\nHowever, previously known dual-chamber and multi-chamber pacemakers do not sense ventricular activity through the atrial lead for a particular interval of time (i.e., the xe2x80x9cpost-ventricular atrial refractory period,xe2x80x9d commonly known as PVARP) subsequent to the delivery of the ventricular stimulation pulse. This refractory period on the atrial channel following ventricular stimulation prevents the atrial channel from mistaking a far-field R-wave for an atrial P-wave. However, detection of the far-field R-wave can be advantageous in verifying that ventricular capture has occurred.\nIt would thus be desirable to provide a system and method for automating the detection of capture on one or both ventricular channels of an implantable multi-chamber stimulation device, with increased accuracy. It would also be desirable to provide a system and method for reducing the negative effect of polarization and noise on capture verification. It would further be desirable to enable the stimulation device to perform ventricular capture verification without requiring dedicated circuitry and/or special sensors.\nThe present invention addresses these needs by providing a system and method for automatically detecting capture of a ventricular chamber in a multi-chamber cardiac stimulation device. Ventricular capture is detected by sensing the far-field R-wave that follows a ventricular stimulation pulse that has successfully captured the ventricle.\nIn one illustrative embodiment of the invention, the device delivers a ventricular stimulation pulse to both the right and left ventricles and then samples a far-field R-wave, resulting from the biventricular evoked response, on the atrial channel during a predetermined far-field interval window.\nIn another illustrative embodiment, the device delivers a ventricular stimulation pulse to one ventricle and then samples a far-field evoked response on a pair of atrial channels during a predetermined far-field interval window.\nIn certain embodiments, capture may be verified by comparing one or more signal characteristics of the far-field R-wave sample, such as peak amplitude, integral, or slope, to the same characteristic(s) of an expected far-field R-wave following an evoked response in the designated ventricle."}
{"text": "The present disclosure generally relates to dental care. More specifically, the present disclosure relates to oral care products for the improved mineralization/remineralization of teeth.\nEnamel and dentin in teeth are primarily composed of calcium phosphate in the form of calcium hydroxyapatite. This material is highly insoluble at neutral salivary pH levels, but tends to dissolve in acidic media. Consequently, when teeth are exposed to acids generated during the bacterial-induced glycolysis of carbohydrates, lesions or demineralized areas are initiated below the surface of intact enamel since the outer rim is more acid resistant. Dental caries begin with these subsurface lesions, which are formed before a cavity is even detectable. If not treated, the surface enamel above such a subsurface lesion will eventually collapse, resulting in the formation of a cavity and subsequent loss of tooth structure."}
{"text": "This application generally relates to communications and, more particularly, to detecting encrypted path finding or routing a message with an address header.\nEncryption of communications is increasing. More and more people, businesses, and governments are encrypting their electronic communications. This encryption provides enhanced security and privacy for these electronic communications.\nEncryption, however, is a problem for communications service providers. Communications service providers need to know the type of data contained within an electronic communication. Some data types receive priority processing, while other data types are queued for later processing. Encryption, however, hides the contents of the communication and often prevents a communications service provider from determining the level of required processing. Because the communications service provider cannot determine the level of required processing, the encrypted communication defaults to lesser priority and/or processing.\nInternet telephony provides an example. Internet telephone calls should be processed to result in a real time, or nearly real time, conversation. If packets are lost, or if packets experience congestion, the quality of the call suffers. Internet telephone calls, then, should receive priority processing. When a communications service provider detects data representing an Internet telephone call, the service provider gives that data priority/special processing to reduce packet loss and to reduce latency effects. Encryption, however, hides the contents of the communication. Encryption prevents the communications service provider from determining whether priority and/or special processing is required. So, even though the communication is an Internet telephone call, encryption causes the communication to default to lesser priority and/or processing. The quality of the call may then suffer from packet loss and congestion.\nThere is, accordingly, a need in the art for improved determination of data types. When parties encrypt their communications, there is a need for determining the type of data contained inside the encrypted communication. There is also a need for identifying a particular kind of encrypted traffic in order to provide prioritized/specialized processing."}
{"text": "Dynamic range is the ratio of intensity of the highest luminance parts of a scene and the lowest luminance parts of a scene. For example, the image projected by a video projection system may have a maximum dynamic range of 300:1.\nThe human visual system is capable of recognizing features in scenes which have very high dynamic ranges. For example, a person can look into the shadows of an unlit garage on a brightly sunlit day and see details of objects in the shadows even though the luminance in adjacent sunlit areas may be thousands of times greater than the luminance in the shadow parts of the scene. To create a realistic rendering of such a scene can require a display having a dynamic range in excess of 1000:1. The term “high dynamic range” means dynamic ranges of 800:1 or more.\nModern digital imaging systems are capable of capturing and recording digital representations of scenes in which the dynamic range of the scene is preserved. Computer imaging systems are capable of synthesizing images having high dynamic ranges. However, current display technology is not capable of rendering images in a manner which faithfully reproduces high dynamic ranges.\nBlackham et al., U.S. Pat. No. 5,978,142 discloses a system for projecting an image onto a screen. The system has first and second light modulators which both modulate light from a light source. Each of the light modulators modulates light from the source at the pixel level. Light modulated by both of the light modulators is projected onto the screen.\nGibbon et al., PCT Application No. PCT/US01/21367 discloses a projection system which includes a pre modulator. The pre modulator controls the amount of light incident on a deformable mirror display device. A separate pre-modulator may be used to darken a selected area (e.g. a quadrant).\nThere exists a need for cost effective displays capable of reproducing a wide range of light intensities in displayed images."}
{"text": "Various hair implantation procedures and equipment associated with the implementation thereof have been employed in the prior art. Follicular unit hair transplantation is the procedure most commonly used. This procedure is however invasive since it involves removal of a strip of the scalp from a donor area and placement of the graft at a recipient area. Such strip removal necessitates shaving and leaves a scar that is difficult to heal, thereby posing restrictions in the client's normal activities over a rather long healing period. The process is rather slow with short daily sessions enabling implantation of a limited number of hair of the order of 500 units, whilst its efficiency is rather low due to a large percentage of up to 20% of follicles being damaged during the removal of the strip and the subsequent division thereof in individual follicles to be implanted into the recipient zone. Eventually since a number of sessions must take place to perform the overall hair transplantation process and since such sessions have to be scheduled with extensively long intervals in between, it becomes uncomfortably long to complete the process. Furthermore, placement of hair follicles is subsequently performed by means of forceps thereby leaving much room for human error that often leads to oddly directioned hair at incorrect angles instead of the desired natural result of evenly angled hair. Moreover such graft extraction and subsequent implanting process is not available for all clients, e.g. very curly hair cannot be extracted, and a preliminary test has to be conducted to determine whether the process is suitable for each particular candidate.\nThe object of the present invention is to overcome and eliminate the abovementioned disadvantages and drawbacks of the prior art. With this scope in mind, the invention proposes a minimally invasive hair transplantation process that is available to all potential candidates without exception and requires neither shaving of the scalp nor removal of grafts or any preliminary testing whatsoever. Instead single hair follicles are extracted from a donor region to be subsequently implanted at a recipient region. The process is fast providing the possibility of an approximate number of 5,000 hair units being extracted in one session.\nThe object of the invention is further to disclose preferred specialized tools for the implementation of the abovementioned direct hair implantation process and in particular to propose alternative embodiments of tool devices as follows:                a. a hair harvesting device adapted to perform cutting of individual hair follicles comprising a tubular cutting head with a conical knife edge cutting surface formed at one end thereof or with conical knife edge cutting surfaces of different diameters formed at two ends thereof.        b. a hair harvesting device comprising a cross like pattern of tubular cutting heads, each one being formed with a conical knife edge cutting surface of different diameter allowing removal of follicular hair units of varying diameters with a single harvesting instrument.        c. a hair implanting device comprising a hollow needle having a longitudinal groove at a frontal portion thereof with an obliquely cut free end adapted to receive a follicular hair unit, a sliding rod being reciprocatingly movable axially along the needle to effect placement of the follicular hair unit disposed within the groove of the frontal end of the needle within a predetermined position of the scalp at the recipient bald area, and a tubular housing adapted to receive said needle and said sliding rod incorporating means of performing the hair implanting operation.        \nIt is a further object of the invention to provide the aforementioned hair implanting device with means of appropriately regulating and fine adjustment of the depth of intrusion of the needle into the scalp to provide for optimum performance of the hair implanting device.\nAnother object of the invention is to disclose an electrically operated mode of the hair transplanting device.\nA final object of the invention, in view of accelerating the hair transplantation process by avoiding the step of the tiresome exchange of instruments, is to provide a combined hair harvesting and hair implanting device operated preferably manually when effecting hair harvesting and preferably electrically when effecting hair implantation, wherein the aforementioned knife edge cutting surface of the hair harvesting device is obliquely cut so as to be adapted to perform the combined role of the follicular hair unit receiving groove of the hair implanting device."}
{"text": "1. Field of the Invention\nThis invention relates generally to wireless and long distance carriers, Internet service providers (ISPs), and information content delivery services/providers and long distance carriers. More particularly, it relates to location services for the wireless industry.\n2. Background of Related Art\nThe Location Interoperability Forum (LIF), the Wireless Application Protocol (WAP) Forum, and 3rd Generation Partnership Project (3GPP) have attempted to define an area trigger via Application Protocol Interfaces (APIs) specific to these groups.\nThe problem with the above solutions is that they are not well defined or do not implement schemes that would permit functionality to support area watching features such as are provided by the present invention."}
{"text": "A middlebox is a network appliance that manipulates internet traffic by optimizing data flow across the network. Middleboxes can be configured as wide area network (“WAN”) optimizers and can be deployed in pairs across two geographically separated locations to optimize data traffic between the two middleboxes. Middleboxes can be connected through a single link or multiple links such as a leased line link and a broadband link. Middleboxes use TCP congestion avoidance algorithms, commonly called “TCP flavors,” to optimize TCP data flows as part of a quality of service (“QoS”) scheme. Common examples of TCP avoidance flavors can include algorithms such as TCP Vegas, TCP Reno, TCP NewReno, TCP Hybla, TCP BIC, and TCP CUBIC, among others. Each TCP congestion avoidance flavor is suited for optimizing data flows originating from or received by particular operating systems, link types, and/or other network characteristics.\nSome TCP flavors improve quality of service across TCP connections by using congestion control and congestion avoidance techniques that sometimes include TCP traffic prioritization. Traffic prioritization, a traffic shaping technique for Quality of Service (QoS), can ensure that more packets from very high priority (P1) traffic are pushed into the network for increased throughput. Thus when there are enough data to be sent from both P1 and non-high priority (non-P1) traffic, the prioritization functionality of QoS can push more packets from P1 traffic into the network pipe thereby delivering better throughput and providing enhanced QoS for P1 traffic. For example, using conventional methods of prioritization, if traffic prioritization module dictates 3:1 ratio for P1 and non-P1 traffic, we could expect approximately 75% of the packets occupying the leased network pipe to be of P1 traffic.\nOne down side of conventional prioritization techniques is when the network link gets congested, since there are more packets from the P1 traffic, the probability of packet drops from P1 traffic is much higher than the probability of packet drops in non-P1 traffic. For example, in one congestion scenario, the probability of next packet drop to occur in P1 traffic would be 0.75 as opposed to non-P1 traffic with a probability of 0.25. Thus, using conventional traffic prioritization techniques, P1 traffic could get multiple packet drops, which in turn forces the TCP connection to reduce its congestion window. Therefore, the higher priority TCP traffic can experience reduced throughput and degraded end-to-end QoS."}
{"text": "It is known in the art relating to oral swabs to apply vacuum force through a porous foam sleeve or through an aperture in the foam sleeve to remove liquid and other material from the mouth of a patient. Such oral swabs clog as the porous surface of the foam sleeve becomes filled or as the aperture in the foam sleeve becomes restricted with debris.\nIn addition, such oral swabs preclude removal of viscus fluid and/or material from oral recesses and recesses between the teeth and gum line. Often times the vacuum force drawn through the foam sleeve collapses the foam into the device and/or onto itself and causes the swab to fail to continue suctioning. Such side hole devices are difficult to manufacture as the foam sleeve must be properly placed on a vacuum supply tube end to ensure flow through the element. In some cases, an aperture in the foam sleeve must be aligned relative to an aperture in the vacuum supply tube end."}
{"text": "To a transmission line configuring a network for optical communication; an optical amplifier provided in this transmission line; or an optical transceiver, an optical transponder, or the like that performs optical communication, an optical power attenuator may be provided to adjust an optical power.\nFurthermore, as the optical power attenuator, known is one of a reflective type provided with a MEMS (microelectromechanical-system) element that can control a reflection angle of a mirror that reflects a light according to an applied voltage.\nThe MEMS element used in the optical power attenuator of the reflective type is provided with the mirror, which is pivotably supported via a beam, and is configured to be able to adjust the reflection angle of the light by the mirror by an electrostatic force arising according to the applied voltage (see for example patent literature 1).\nFurthermore, as illustrated in FIG. 5, a conventional optical power attenuator 5 is provided with a MEMS element 10 configured as above, a capillary member 22 provided to one end of a two-core optical fiber 20 for inputting/outputting a light, and a lens 30 disposed between the capillary member 22 and the MEMS element 10.\nFurthermore, the lens 30 is configured as a collimating lens; it converts an incident light from an IN-side optical fiber 20i among the two-core optical fiber 20 into a collimated light (parallel light), causes this to become incident to the MEMS element 10, and causes a reflected light from the MEMS element 10 to become incident to an OUT-side optical fiber 20o in a path different from an incoming path.\nBecause of this, the incident light from the IN-side optical fiber 20i is attenuated according to a reflection angle of the light by a mirror in the MEMS element 10, becomes incident to the OUT-side optical fiber 20o, and is transmitted to another optical device via the OUT-side optical fiber 20o. "}
{"text": "1. Field of the Invention\nThe invention relates in general to a touch interface and an operating method thereof, and more particularly to a touch control electronic device and an operating method thereof.\n2. Description of the Related Art\nConventionally, communication interfaces between human and computers or machines are mainly based on keyboards and mice. With the trend of the technological product toward the friendlier man-machine interface, the applications of the touch interfaces become more and more popularized. Recently, the flourishing development of the highly technological industry brings the development of the information and consumer products. The user is eager for an operation interface that can be easily operated so that the touch interface has become the design stream of the product.\nIn the field of the mobile technology, handheld electronic devices, such as a personal digital assistant (PDA), a smart phone, a portable video game, a portable multimedia player and a portable navigation device, also adopt touch panels to satisfy the miniaturized design trend. In addition, the touch inducing manipulation method of the recently developed touch wheel has overcome the conventional mechanical manipulation method, in which typical keys and a scroll wheel are equipped. The touch wheel may directly sense the circular movement of the user's finger. The user only has to put his/her finger on the touch wheel and to slide, rotate or touch the peripheral and middle keys of the wheel to select the items on the screen, to adjust the volume, or to click the songs so that the handheld electronic device can be easily operated.\nThus, it is a new direction for the future product development to effectively apply the touch interface to create the new product using the method so that the user can operate the touch interface more conveniently and instinctively and the man-machine interaction effectiveness of the user can be enhanced."}
{"text": "In the fabrication of ultra-large scale-integration (ULSI) circuits it has been very common to utilize vertical stacking, or vertical integration, of metal wiring circuits to form multilevel interconnection. Multilevel fabrication process has become an efficient way to increase circuit performance and increase the functional complexity of the circuits. One drawback of multilevel interconnection is the loss of topological planarity resulting from various photolithographic and etching processes. The various integrated circuit fabrication processes invariably produce nonplanar surface, or nonplanar topography, on the wafer, from which semiconductor devices are fabricated. During the multilevel metallization of VLSI or ULSI devices, the multiplicity of layers of nonplanar surfaces further add together to cause even more serious topography problems. For example, the conductive or insulative properties of the various deposited films can be degraded on the area of the film layers across the step height. Those films in high topography areas can be easily broken during heat, electrical current, or mechanical stress steps, resulting in the pattern areas becoming discontinuous. Such discontinuity can cause failures in the device to perform certain intended functions. Furthermore, a nonplanar surface cannot be precisely focused during the photolithography process, because the depth of focus of the conventional photolithographic stepper will be deviated by different step heights of the wafer. Such an out-of-focus problem is more profound with respect to device features of very small sizes.\nTo alleviate these problems, the wafer is planarized at various stages in the fabrication process to minimize non-planar topography and thus its adverse effects. Such planarization is typically implemented in the dielectric layers. However, it is also possible to implement the planarization process in the conductor layer. More recently, chemical-mechanical polishing (CMP) processes have become very well received to planarize the wafer surface in preparation for further device fabrication. The CMP process mainly involves the step of holding a semiconductor wafer against a rotating polishing pad surface wetted by a polishing slurry, which typically comprises an acidic or basic etching solution in combination with alumina or silica particles. On the one hand, the liquid portion of the slurry chemically removes, loosens, or modifies the composition of the material on the wafer which is to be removed. On the other hand, the particle portion of the slurry, in combination of the rotating polishing pad, physically removes the chemical modified material from the wafer. Thus, the name \"chemical-mechanical polishing\" was obtained.\nOne of the most commonly used sacrificial materials in the chemical-mechanical polishing process is a solution-type silicon dioxide, which is commonly referred to as the spin-on-glass (SOG). The SiO.sub.2 -based SOG is initially formed as a low viscosity solution which can be coated onto the nonplanar surface to quickly fill the recessed areas by a conventional spin coating technique. After the SOG coating, the coated layer is hard-baked to remove the solvent contained therein and turn the SOG layer into a hardened layer. Because of its high electrical resistance, the solidified SOG layer on top of the integrated circuit structure (i.e., the metal layer) must be etched back, typically by a chemical-mechanical polishing procedure using an abrasive slurry containing hydrogen fluoride in the chemical portion of the polishing procedure. If the SOG layer on the top surface of the wafer is not completely removed, it can generate the so-called vias poisoning, causing the vias to have a very high electrical resistance and adversely affect the interlayer conduction.\nOne of the commonly observed problems in using SOG during the planarization process is the so-called dishing effect. The etching process is not a selective process and it will remove all the affected material on the wafer surface. During the chemical-mechanical polishing of the SOG layer, the portion of the SOG layer inside the trenches can be affected by the etching solution being dished out from the mechanical polisher, resulting another type of un-planarized top surface. Many times, the dishing effect could cause even more serious nonplanar surface on the wafer. One way to ameliorate the dishing effect in the etch back process is to use a etch-back photoresist. However, this would require an extra photolithography process using an extra photo mask to remove the sacrificial material above the patterned structure. This can add substantially to the total manufacturing cost.\nAnother way to avoid such dishing effect is to use a non-silicon-based polymeric material, such as polyimide, which exhibits excellent chemical and electrical resistance, as the sacrificial material for chemical-mechanical polishing processes. The polymer-based sacrificial material, however, will not be removed by the conventional chemical-mechanical polishing slurry designed for SOG. Thus, the IC manufacturers must stock two different types of chemical-mechanical polishing slurries. And the use of a different chemical-mechanical polishing slurry may cause material compatibility problems and other handling concerns. In a co-pending application, it was demonstrated that excellent result can be obtained by treating the polymeric material so that it becomes removable by hydrofluoric acid. However, in semiconductor devices which contain relatively wide trenches, the polymeric sacrificial layer will be warped into the trench bottom, and the dishing effect can again become significant.\nU.S. Pat. No. 4,944,836 disclosed a chemical-mechanical polishing method for producing coplanar metal-insulator films on a substrate. It taught the conventional approach of using silicon dioxide based slurry to provide a planarized surface by chemical-mechanical polishing. However, when the layer thickness of the semiconductor device becomes increasing thin, and the depth of the trench becomes correspondingly shallow, the dishing effect caused by the exposure of the silicon dioxide in the trench area becomes relatively more profound.\nU.S. Pat. No. 5,169,491 disclosed a method for planarizing SiO.sub.2 -containing dielectric in semiconductor wafer processing. It involved the steps of (1) providing a layer of undoped SiO.sub.2 atop a wafer; (b) depositing a layer of borophosphosilicate glass (BPSG) atop the layer of undoped SiO.sub.2 ; and (c) chemical mechanical polishing the borophosphosilicate glass selectively relative to the underlying layer of undoped SiO.sub.2 and using the underlying layer of undoped SiO.sub.2 as an etching stop.\nU.S. Pat. No. 5,314,843 disclosed a method for planarizing a semiconductor wafer using the chemical-mechanical polishing process. The method included the steps of (1) masking the semiconductor wafer surface layer to define first and second laterally adjacent portions; (2) altering only the first portion of the surface layer of material to polish at a different removal rate in the chemical-mechanical polishing process than the second portion of the surface layer; and (3) polishing the surface in the chemical-mechanical polishing process. As discussed before, the method disclosed in the '843 patent may minimize the dishing effect, but it requires an additional expensive photolithography step."}
{"text": "Modern vehicles (e.g., airplanes, boats, trains, cars, trucks, etc.) can include a vehicle event recorder in order to better understand the timeline of an anomalous event (e.g., an accident). A vehicle event recorder typically includes a set of sensors, e.g., video recorders, audio recorders, accelerometers, gyroscopes, vehicle state sensors, GPS (global positioning system), etc., that report data, which is used to determine the occurrence of an anomalous event. If an anomalous event is detected, the sensor data related to the event is stored for later review. A vehicle event recorder for cars and trucks (e.g., vehicles that operate on public roads) can include road map data comprising location-specific legal information (e.g., speed limit information, stop sign information, traffic light information, yield sign information, etc.). Location-specific legal information can be used to identify an anomalous event in the case of the vehicle acting against the law (e.g., traveling in excess of the speed limit, rolling through a stop sign, etc.). If there is an error in the legal information, anomalous events can be incorrectly identified, possibly leading to unnecessary expense as the event is processed, stored, and/or transmitted."}
{"text": "1. Field of the Invention\nThis invention relates to aroma therapy delivery systems, and more particularly to aroma therapy delivery systems incorporating a shaving unit as the delivery vehicle.\n2. Description of the Background Art\nVarious conventional wet shavers disclose the general concept of applying scented substances to a user's skin as lubricating agents during the shaving process.\nFor example, U.S. Pat. No. 4,850,107 issued to Valliades et al. on Jul. 25, 1989 discloses a razor assembly with means for intermittently distributing a thin fluid film beneath the bottom of the blade while shaving. The '107 patent discloses an open recessed seating area on the upper, angled portion of the razor shaft for housing a sponge. By applying pressure to the sponge, a thin fluid is released through channels onto the face. The '107 patent teaches that the thin fluid, which may be a scented fluid, may be released on demand to moisten the skin of the user while shaving.\nU.S. Pat. No. 4,875,287 issued to Creasy et al. on Oct. 24, 1989 discloses a razor head having a coated surface or substrate which provides, inter alia, a lubricant to the user's face. The '287 patent discloses that the additional materials can be incorporated into the polymer blends such as fragrances.\nU.S. Pat. No. 5,134,775 issued to Althaus et al. on Aug. 4, 1992 discloses a shaver head for a wet shaver comprising a device for receiving a liquid shaving preparation which is dispensed during shaving. The '775 patent discloses that the liquid shaving preparation can be perfumed.\nU.S. Pat. No. 5,121,541 issued to Patrakis on Jun. 16, 1992 discloses an electric razor comprising a misting mechanism for misting a lubricating agent, such as water, cologne or a beard softener, onto the user's skin while shaving.\nWhile these publications appear to disclose the general concept of using scented lotions or lubricating agents to be applied to a user's face while shaving, they fail to teach an aroma therapy delivery system or aroma therapy delivery systems incorporating a shaving unit as the delivery vehicle which dispense an aroma therapy.\nAccordingly, it is the primary object of the present invention to provide a shaver assembly unit which dispenses an aromatic therapy.\nIt is a further object of the invention to provide a shaver assembly unit that can be easily filled with an aromatic agent."}
{"text": "1. Field of Art\nThe disclosure generally relates to the field of data collection and data mining in computing devices.\n2. Description of Art\nAs mobile computing technology advances, more and more applications become available for mobile computing devices. As a result, users use the mobile computing devices to perform more activities. In addition, mobile computing devices are also equipped with increasing number of sensors such as Global Positioning System (GPS) receivers, accelerometers, and proximity sensors. These sensors, coupled with the increased number of applications, give users access to even more information than has previously been available.\nThe information from the different sources has inherently relationships (e.g., the GPS receiver tracks the current geographic location while the calendar application provides the user's scheduling information). Currently, there is no solution to automatically collect and analyze the user information available on multiple sources of the mobile computing device to discover the inherent relationships and to provide the user with summaries of the collected user information and the relationships. Accordingly, there is lacking in the art, inter alia, techniques for collecting and analyzing user activities on a mobile computing device."}
{"text": "Recent fuel developments have resulted in a number of aqueous fuel emulsions comprised essentially of a carbon based fuel, water, and various additives such as lubricants, emulsifiers, surfactants, corrosion inhibitors, cetane improvers, and the like. These aqueous fuel emulsions may play a key role in finding a cost-effective way for internal combustion engines including, but not limited to, compression ignition engines (i.e. diesel engines) to achieve the reduction in emissions below the mandated levels without significant modifications to the engines, fuel systems, or existing fuel delivery infrastructure.\nAdvantageously, aqueous fuel emulsions tend to reduce or inhibit the formation of nitrogen oxides (NOx) and particulates (i.e. combination of soot and hydrocarbons) by altering the way the fuel is burned in the engine. Specifically, the fuel emulsions are burned at somewhat lower temperatures than a conventional fuels due to the presence of water. This, coupled with the realization that at higher peak combustion temperatures, more NOx are typically produced in the engine exhaust, one can readily understand the advantage of using aqueous fuel emulsions.\nA major concern of aqueous fuel emulsions or water blend fuels, however, is the stability of the fuel. As is well known in the art, the constituent parts of such aqueous fuel emulsions have a tendency to separate over time. Blending of the fuel emulsions in a manner to achieve long-term stability is essential if such fuels are to be commercially successful. The problems associated with fuel emulsion separation are very severe inasmuch as most engine operating characteristics are adjusted for a prescribed fuel composition. Where the fuel emulsion composition has changed due to ingredient separation, the engine performance is markedly diminished.\nSeveral related art references have disclosed various devices or techniques for producing or blending a fuel emulsion for internal combustion engines. For example, U.S. Pat. No. 5,535,708 (Valentine) discloses a process for forming an emulsion of an aqueous urea solution in diesel fuel and combusting the same for the purposes of reducing NOx emissions from diesel engines. See also U.S. Pat. No. 4,938,606 (Kunz) discloses an apparatus for producing an emulsion for internal combustion engines that employs an oil line, a water line, a dosing apparatus and various mixing and storage chambers. Another related art process and system for blending a fuel emulsion is disclosed in U.S. Pat. No. 5,298,230 (Argabright) which discloses a specialized process for blending a fuel emulsification system useful for the reduction of NOx in a gas turbine.\nThe present invention addresses the aforementioned problems associated with separation of aqueous fuel emulsions by providing a blending system and method that enhances the long term stability of such emulsions."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetoresistive element using a method of reducing magnetic field inverting magnetization thereinafter, referred to as switching field of a magnetic film, a memory element having the magnetoresistive element, and a memory using the memory element.\n2. Related Background Art\nIn recent years, semiconductor memories as solid-state memories are adopted in many information devices, and are of various types such as a DRAM, FeRAM, and flash EEPROM. The characteristics of the semiconductor memories have merits and demerits. There is no memory which satisfies all specifications required by current information devices. For example, the DRAM achieves high recording density and large rewritable count, but is volatile and loses its information upon power-off. The flash EEPROM is nonvolatile, but takes a long erase time and is not suitable for high-speed information processing.\nUnder the present circumstances of semiconductor memories, a magnetic memory (MRAM: Magnetic Random Access Memory) using a magnetoresistive element is promising as a memory which satisfies all specifications required by many information devices in terms of the recording time, read time, recording density, rewritable count, power consumption, and the like. In particular, an MRAM using a spin-dependent tunneling magnetoresistive (TMR) effect is advantageous in high-density recording or high-speed read because a large read signal can be obtained. Recent research reports verify the feasibility of MRAMs.\nThe basic structure of a magnetoresistive film used as an MRAM element is a sandwich structure in which magnetic layers are formed adjacent to each other via a nonmagnetic layer. Known examples of the material of the nonmagnetic film are Cu and Al2O3. A magnetoresistive film using a conductor such as Cu in a nonmagnetic layer is called a GMR film (Giant MagnetoResistive film). A magnetoresistive film using an insulator such as Al2O3 is called a spin-dependent TMR film (Tunneling MagnetoResistive film). In general, the TMR film exhibits a larger magnetoresistance effect than the GMR film.\nWhen the magnetization directions of two magnetic layers are parallel to each other, as shown in FIG. 13A, the resistance of the magnetoresistive film is relatively low. When these magnetization directions are antiparallel, as shown in FIG. 13B, the resistance is relatively high. One of the magnetic layers is formed as a recording layer, and the other layer is as a read layer. Information can be read out by utilizing the above property. For example, a magnetic layer 13 on a nonmagnetic layer 12 is formed as a recording layer, and a magnetic layer 14 below the nonmagnetic layer 12 is as a read layer. The rightward magnetization direction of the recording layer is defined as xe2x80x9c1xe2x80x9d, and the leftward direction is as xe2x80x9c0xe2x80x9d. If the magnetization directions of the two magnetic layers are rightward, as shown in FIG. 14A, the resistance of the magnetoresistive film is relatively low. If the magnetization direction of the read layer is rightward and that of the recording layer is leftward, as shown in FIG. 14B, the resistance is relatively high. If the magnetization direction of the read layer is leftward and that of the recording layer is rightward, as shown in FIG. 14C, the resistance is relatively high. If the magnetization directions of the two magnetic layers are leftward, as shown in FIG. 14D, the resistance is relatively low. That is, when the magnetization direction of the read layer is pinned rightward, xe2x80x9c0xe2x80x9d is recorded in the recording layer for a high resistance, and xe2x80x9c1xe2x80x9d is recorded for a low resistance. Alternatively, when the magnetization direction of the read layer is pinned leftward, xe2x80x9c1xe2x80x9d is recorded in the recording layer for a high resistance, and xe2x80x9c0xe2x80x9d is recorded for a low resistance.\nAs the element is downsized for a higher recording density of an MRAM, the MRAM using an in-plane magnetization film becomes more difficult to hold information under the influence of a demagnetizing field or magnetization curling at the end face. To avoid this problem, for example, a magnetic layer is formed into a rectangle. This method cannot downsize the element, so an increase in recording density cannot be expected. U.S. Pat. No. 6,219,275 has proposed the use of a perpendicular magnetization film to avoid the above problem. According to this method, the magnetizing field does not increase even with a smaller element size. A smaller-size magnetoresistive film can be realized, compared to an MRAM using an in-plane magnetization film. Similar to a magnetoresistive film using an in-plane magnetization film, a magnetoresistive film using a perpendicular magnetization film exhibits a relatively low resistance when the magnetization directions of two magnetic layers are parallel to each other, and a relatively high resistance when these magnetization directions are antiparallel. As shown in FIGS. 15A to 15D, a magnetic layer 23 on a nonmagnetic layer 22 is formed as a recording layer, and a magnetic layer 21 below the nonmagnetic layer 22 is as a read layer. The upward magnetization direction of the recording layer is defined as xe2x80x9c1xe2x80x9d, and the downward direction is as xe2x80x9c0xe2x80x9d. As FIGS. 14A to 14D showed, it can compose as a memory element.\nMain examples of the perpendicular magnetization film are an alloy film or artificial lattice film made of at least one element selected from the group consisting of rear-earth metals such as Gd, Dy, and Tb and at least one element selected from the group consisting of transition metals such as Co, Fe, and Ni, an artificial lattice film made of a transition metal and noble metal such as Co/Pt, and an alloy film having crystallomagnetic anisotropy in a direction perpendicular to the film surface, such as CoCr. In general, the switching field of a perpendicular magnetization film is larger than that having longitudinal magnetic anisotropy by a transition metal. For example, the switching field of a permalloy as an in-plane magnetization film is about several hundred A/m. The switching field of a Co/Pt artificial lattice film as a perpendicular magnetization film is as very high as about several ten kA/m. An alloy film of a rear-earth metal and transition metal exhibits different apparent magnetization intensities depending on the film composition because the sub-lattice magnetization of the rear-earth metal and that of the transition metal orient antiparallel to each other. Hence, the switching field of this alloy film changes depending on the composition. A GdFe alloy film shows a relatively small switching field among alloy films of rear-earth metals and transition metals. In general, the GdFe alloy film has a switching field of about several thousand A/m around the critical composition at which the squareness ratio of the magnetization curve starts decreasing from 1.\nWhen a sensor, memory, or the like is formed from a magnetoresistive film using a perpendicular magnetization film, the sensor, memory, or the like cannot operate unless a large magnetic field is applied owing to the above-described reason. For example, in the sensor, a stray field must be concentrated on the magnetic layer of the magnetoresistive film. In the memory, a large magnetic field must be generated. A magnetic field applied to a memory is generally generated by supplying a current through a conductor. Especially in a memory used in a portable terminal, supply of a large current is undesirably flowed under restrictions on the power supply capacity. Thus, a conductor for generating a magnetic field must be wound around a memory element formed from a magnetoresistive film. This measure complicates a structure or electrical circuit around the magnetoresistive film, and is difficult to form. This results in low yield and very high cost.\nThe present invention has been made in consideration of the above situation, and has as its object to provide a magnetoresistive film which reduces the switching field of a perpendicular magnetization film and is easy to form without decreasing the yield or greatly increasing the cost, and a memory requiring only small power consumption.\nThe above object is achieved by a magnetoresistive film comprising a nonmagnetic film, and a structure in which magnetic films are formed on two sides of the nonmagnetic film, wherein at least one of the magnetic films includes a perpendicular magnetization film, and a magnetic film whose easy axis of magnetization is inclined from a direction perpendicular to a film surface is formed at a position where the magnetic film contacts the perpendicular magnetization film but does not contact the nonmagnetic film.\nThe above object is also achieved by a memory having a memory element with the magnetoresistive film, comprising means for applying a magnetic field to the magnetoresistive film in a direction perpendicular to a film surface, and means for detecting a resistance of the magnetoresistive film.\nThe above object is also achieved by the memory wherein a plurality of magnetoresistive films are arranged, and the memory further comprises means for selectively recording information on a desired magnetoresistive film, and means for selectively reading out information recorded on a desired magnetoresistive film."}
{"text": "The invention relates generally to molding forms and, more particularly, to an apparatus for constructing a molding form having an accurately located reference surface that establishes a border for the molded structure.\nMolded structures formed along the surface of the ground, such as concrete decks, typically are constructed by preparing a firm bed upon which the deck is to rest. Border defining forms are provided for confining the uncured concrete mixture to a desired border contour. A concrete mixture is poured over the prepared bed to fill the form and the surface of the poured concrete troweled. After the concrete cures to a firm state, the forms are dismantled.\nTwo types of molding form structures have been used in constructing concrete decks; structures of wood form boards and stakes fastened together by nails, and structures of specially formed metal form boards and stakes held together by rigid metal wedges. The wood form structures are constructed by driving stakes into the ground at intervals along a stake line that follows the desired edge contour of the deck. Then, flexible form boards of wood, commonly referred to as bender boards, are placed against the stakes and fastened thereto by nailing. In most cases, the stakes are displaced by the force of nailing, making it virtually impossible to locate the reference surface of the bender board at the reference line for the border of the deck. In addition to not being able to locate the reference surface precisely, the manner of construction of such wood form structures causes rapid deterioration of the materials used. Repeated nailing weakens the wood bender boards and they frequently break during the assembling and dismantling of the forms. The low number of repeated uses of the wood bender boards and stakes and the consumption of nails is an expense of construction that is becoming increasingly significant with the rising costs of construction materials.\nThe metal type molding form structure includes a metal form board having a pair of apertured appendages extending perpendicularly from one side of a planar member, the opposite side of the planar member forming the reference surface. A support stake is inserted through apertures of the appendages and a separate rigid metal wedge is driven between the stake and the facing side of the metal form board. Driving the wedge forces the metal form board away from the support stake to cause the stake to bear forcefully against the wall of the apertures whereby the metal form board is held firmly in place relative to the support stake. A substantial force is required to drive the rigid wedge between the stake and the facing side of the form board and, frequently, the wedge galls the side and/or support stake. After many uses of the metal form boards and support stakes, galling often becomes so extensive that the metal form boards and/or support stakes are deteriorated beyond use. Even before the metal form boards and/or support stakes are deteriorated beyond use, galling may deform them to such an extent that accurate location of the reference surface that defines the border of the concrete deck is prevented. Consequently, the metal type molding form structures do not entirely alleviate the shortcomings of the previously described wood type. In addition, it is more costly to manufacture and use the metal form structures than their wood counterparts.\nAccordingly, considerable advantage is to be gained by the use of inexpensive molding form components that are not easily damaged by repeated usage. Additional advantages are to be gained by the use of molding form components that permit commonly available bender boards to be used repeatedly."}
{"text": "1. Field of the Invention\nThe invention relates to a position sensing device and a motor using the same.\n2. Description of the Related Art\nSensing devices such as are widely used in motors to detect positions of stators. However, there are several problems with conventional sensing devices: structure of the sensing devices is complex, and connection and fixation between the sensing devices and electric components such as power leads, PCB boards and so on are not firm enough."}
{"text": "Technical Field\nThe present disclosure relates to memory devices such as a semiconductor memory, and method for testing reliability of memory devices.\nRelated Art\nIn general reliability tests on semiconductor memories, a testing device is used to write and read data to and from all regions in a memory array with a known test pattern, and the data written to the memory array by the testing device (expected value) is compared with the data read from memory array by the testing device, so as to check the reliability of the memory array.\nIn pre-shipment inspection of semiconductor memories, in order to reduce a testing cost, reliability test is generally performed concurrently on multiple semiconductor memories, by connecting multiple semiconductor memories to one testing device and writing and reading data to and from the multiple semiconductor memories with a common test pattern.\nWith semiconductor memories provided with a pseudo-random number generator for improving security, random number values of pseudo-random numbers are predictable, and thus multiple semiconductor memories can generate an identical pseudo-random number by using a common algorithm. Pre-shipment inspection can therefore be conducted concurrently on multiple semiconductor memories with one testing device, in the same way as on general semiconductor memories.\nRandom number generators are cryptographic technology employed for a wide variety of uses in many security systems.\nRandom numbers generated by random number generators are used for, for example, key information in a cryptographic algorithm, or authentication codes for mutual authentication between devices, and are closely related to the security strength of a system and thus highly confidential information.\nRandom numbers generated by random number generators therefore need to be highly random. At shipment of semiconductor devices provided with a random number generator, a random number test is normally performed to evaluate whether a random number generator generates random numbers that meet a required level.\nJP2005-517998A and WO2005/124537A describe a technique to evaluate whether the frequencies of appearance of “0” and “1” in random numbers generated by a random number generator are within an allowable range."}
{"text": "Wireless and voice-over-internet protocol (VoIP) communications are subject to frequent loss of packets as a result of adverse connection conditions. Such lost packets result in clicks and pops or other artifacts being present in the output voice signal at the receiving end of the connection. This degrades the perceived speech quality at the receiving end and may render the speech unrecognizable if the packet loss rate is sufficiently high.\nBroadly speaking, two approaches are taken to combat the problem of lost packets. The first approach is the use of transmitter-based recovery techniques. Such techniques include retransmission of lost packets, interleaving the contents of several packets to disperse the effect of packet loss, and addition of error correction coding bits to the transmitted packets such that lost packets can be reconstructed at the receiver. In order to limit the increased bandwidth requirements and delays inherent in these techniques, they are often employed such that packet loss can be recovered if the packet loss rate is low, but not all packet loss can be recovered if the packet loss rate is high. Additionally, some transmitters may not have the capacity to implement transmitter-based recovery techniques.\nThe second approach taken to combating the problem of lost packets is the use of receiver-based concealment techniques. Such techniques are generally used in addition to transmitter-based recovery techniques to conceal any remaining losses left after the transmitter-based recovery techniques have been employed. Additionally, they may be used in isolation if the transmitter is incapable of implementing transmitter-based recovery techniques. Low complexity receiver-based concealment techniques such as filling in a lost packet with silence, noise, or a repetition of the previous packet are used, but result in a poor quality output voice signal. Regeneration based schemes such as model-based recovery (in which speech on either side of the lost packet is modeled to generate speech for the lost packet) produce a very high quality output voice signal but are highly complex, consume high levels of power and are expensive to implement. In practical situations interpolation-based techniques are preferred. These techniques generate a replacement packet by interpolating parameters from the packets on one or both sides of the lost packet. These techniques are relatively simple to implement and produce an output voice signal of reasonably high quality.\nPitch based waveform substitution is a preferred interpolation-based packet loss recovery technique. The pitch period of the voiced packets on one or both sides of the lost packet is estimated. A waveform of the estimated pitch period is then repeated and used as a substitute for the lost packet. This technique is effective because voice signals appear to be composed of a repeating segment when viewed over short time intervals. Consequently, the pitch period of the lost voice packet will normally be substantially the same as the pitch period of the voice packets on either side of the lost packet.\nMany methods are used to estimate the pitch period of a voice signal. Generally speaking, these methods include use of a normalized cross-correlation (NCC) method. Such a method can be expressed mathematically as:\n                              N          ⁢                                          ⁢          C          ⁢                                          ⁢                                    C              t                        ⁡                          (              τ              )                                      =                                            ∑                              n                =                                                      -                    N                                    /                  2                                                                              (                                      N                    /                    2                                    )                                -                1                                      ⁢                                          x                ⁡                                  [                                      t                    +                    n                                    ]                                            ⁢                              x                ⁡                                  [                                      t                    +                    n                    -                    τ                                    ]                                                                                                        ∑                                  n                  =                                                            -                      N                                        /                    2                                                                                        (                                          N                      /                      2                                        )                                    -                  1                                            ⁢                                                                    x                    2                                    ⁡                                      [                                          t                      +                      n                                        ]                                                  ⁢                                                      ∑                                          n                      =                                                                        -                          N                                                /                        2                                                                                                            (                                                  N                          /                          2                                                )                                            -                      1                                                        ⁢                                                            x                      2                                        ⁡                                          [                                              t                        +                        n                        -                        τ                                            ]                                                                                                                              (                  equation          ⁢                                          ⁢          1                )            where x is the amplitude of the voice signal and t is time. The equation represents a correlation between two segments of the voice signal which are separated by a time τ. Each of the two segments is split up into N samples. The nth sample of the first segment is correlated against the respective nth sample of the other segment.\nThis equation essentially takes a first segment of a signal (marked A on FIG. 1) and correlates it with each of a number of further segments of the signal (for ease of illustration only three, marked B, C and D, are shown on FIG. 1). Each of these further segments lags the first segment along the time axis by a lag value (τ1 for segment B, ρ2 for segment C). The calculation is carried out over a range of lag values within which the pitch period of the voice signal is expected to be found. The term on the bottom of the fraction in equation 1 is a normalizing factor. The lag value τNCC that maximizes the NCC function represents the time interval between the segment A and the segment with which it is most highly correlated (segment D on FIG. 1). This lag value τNCC is taken to be the pitch period of the signal.\nCalculation of the normalized cross-correlation accounts for over 90% of the algorithmic complexity in typical pitch based waveform substitution techniques. Although the complexity level of the calculation is low, it is significant for low-power platforms such as Bluetooth. In order to correctly determine the pitch period of a voice signal, a wide pre-defined pitch period range (range of lag values) is usually used, for example from 2 ms (for a person with a high voice) to 20 ms (for a person with a low voice). For most pitch determination algorithms, the wider the pitch period range used, the higher the computational complexity.\nOne way to reduce the computational complexity is to reduce the number of calculations that the algorithm computes. U.S. patent application Ser. No. 10/394,118 proposes to reduce the number of calculations by dynamically adapting the time interval between successive segments that are correlated with the first segment. (In the illustration of FIG. 1, the time interval between successive segments B and C is τ2-τ1.) If the correlation decreases, then the time interval to the next segment to be correlated is increased. Conversely, if the correlation increases, then the time interval to the next segment is decreased. This method evaluates the correlation over the same range of pitch periods (for example from 2 ms-20 ms) as methods in which the time interval between successive segments is constant, but advantageously this method is less computationally complex because it carries out fewer calculations by skipping over segments that it considers unlikely to lag the first segment by the pitch period. However, this method is sensitive to local pitch errors. For example, if an error leads to the correlation decreasing just before the pitch period lag value is computed, then the time interval to the next segment may be increased resulting in the algorithm skipping over the pitch period lag value. The accuracy of the estimated pitch period may suffer as a result. Additionally, this method may have difficulty handling voice signals with rapid local pitch variations.\nA further problem with pitch based waveform substitution techniques is that they are prone to pitch doubling and pitch halving errors. Pitch halving occurs when the pitch period is determined to be about double its actual length. This may occur, for example with the method described by U.S. Ser. No. 10/394,118 if the peak best correlated with the peak in the first segment were to be skipped over.\nPitch doubling occurs when the pitch period is determined to be about half its actual length. This may happen in the following situation. Voice signals often have two similar peaks per pitch period that are highly correlated with each other. For example, on FIG. 1 the peaks marked 1 and 2 are highly correlated. These could be mistaken for being the same feature present in consecutive pitch periods and hence the time interval between them could be computed to be the estimated pitch period of the signal. Pitch doubling is particularly problematic for packet loss concealment applications because the replacement signal used for the lost packet will be at a non-integer multiple of the pitch period of the lost packet.\nTechniques for reducing pitch doubling and pitch halving errors have been proposed, for example frequency domain and statistical techniques and post processing techniques. However these techniques incur additional computational complexity and cost.\nThere is thus a need for an improved method of estimating the pitch period of a signal that reduces the computational complexity associated with the estimation, and that additionally reduces susceptibility to pitch doubling and pitch halving errors without incurring extra algorithmic complexity."}
{"text": "1. Field of the Invention\nThe invention relates to a combined power station installation with a gas turbine and a steam turbine, in which the exhaust gases from the gas turbine give up their residual heat to the steam turbine via the working medium flowing in a waste heat boiler, whereby the waste heat boiler consists essentially of an economizer, an evaporator and a superheater and whereby at least one cooling air cooler is provided which is designed as a forced circulation steam generator and is connected on the water side to the economizer of the waste heat boiler.\n2. Discussion of Background\nGas turbines of the modern generation and the higher power class operate with very high turbine inlet temperatures, which makes cooling of the combustion chambers, the rotors and the blading unavoidable. For this purpose, highly compressed air is generally extracted at the compressor outlet and, if appropriate, from a lower pressure stage. Because a very high proportion of the compressed air is consumed for the currently conventional premixed combustion, there remains--on the one hand--only a minimum of cooling air for cooling purposes. On the other hand, this air intended for cooling is already very hot because of the compression so that it is desirable that it should be previously cooled. Cooling by means of water injection (\"gas quenching\") is known for this purpose; in this method, however, the high-quality heat of the cooling air, whose proportion can amount to as much as 20 MW in current machines, is only partially utilized. In consequence, the use of forced circulation steam generators as coolers for recooling seems appropriate, particularly if the gas turbine operates in a combined gas/steam turbine process with waste heat steam generation.\nSuch a once-through steam generator for cooling highly compressed air of the type mentioned at the beginning is known, in association with a combined gas/steam turbine process, from EP-A-709 561. In this specification, a partial flow of the boiler feed water is extracted either upstream or downstream of the economizer and, after further preheating, evaporation and superheating in the cooler, is fed back into the high pressure superheater of the waste heat boiler. This boiler is designed as a circulating system boiler with drums. In order to avoid the penetration of moisture or water into the steam turbine when the cooler is run wet, the heated water or wet steam is fed into a blow-down tank until the cooler is dry or until defined conditions are stably present at the cooler outlet, for example hot steam with a few degrees Kelvin superheat or wet steam with a humidity of a few percent. In addition to the water losses, this has the consequent disadvantage of a corresponding monitoring and control system."}
{"text": "In recent years, as a battery anticipated to have small size, light weight, and high capacity, a non-aqueous electrolytic solution-based secondary battery such as a lithium ion battery has been proposed and put into practical use.\nThe lithium ion battery is constituted of a cathode and an anode which allow the reversible insertion and removal of lithium ions, and a non-aqueous electrolyte.\nRegarding an anode material for lithium ion batteries, as an anode active material, generally, a lithium-containing metal oxide allowing the reversible insertion and removal of lithium ions such as a carbon-based material or lithium titanate (Li4Ti5O12) is used.\nOn the other hand, regarding a cathode material for lithium ion batteries, as a cathode active material, generally, a lithium-containing metal oxide allowing the reversible insertion and removal of lithium ions such as lithium iron phosphate (LiFePO4) or an electrode material mixture is used. In addition, the cathode in the lithium ion battery is formed by applying the electrode material mixture to the surface of a metal foil called a current collector.\nCompared with secondary batteries of the related art such as lead batteries, nickel-cadmium batteries, and nickel-hydrogen batteries, lithium ion batteries have a lighter weight, a smaller size, and higher energy, and thus are used not only as small-size power supplies but also as large-size stationary emergency power supplies in portable electronic devices such as mobile phones and notebook personal computers.\nIn addition, recently, studies have been underway regarding the use of lithium ion batteries as high-output power supplies for plug-in hybrid vehicles, hybrid vehicles, and electric power tools, and batteries used as the high-output power supplies are required to have high-speed charge and discharge characteristics.\nHowever, electrode materials including an electrode active material, for example, a lithium phosphate compound allowing the reversible insertion and removal of lithium ions have a problem of low electron conductivity. Therefore, as electrode materials having increased electron conductivity, an electrode material in which particle surfaces of an electrode active material are uniformly coated with a chemically-deposited carbonaceous material, and the current density of the electrode active material is improved (Japanese Laid-open Patent Publication No. 2001-15111), an electrode material including a carbon black complex obtained by conjugating fibrous carbon and carbon black and olivine-type lithium iron phosphate (Japanese Laid-open Patent Publication No. 2011-108522), and the like have been proposed."}
{"text": "As electronic components continue to advance and are made more powerful, they tend to produce more and more undesirable heat which is preferably removed. This has created a growing need for higher capacity cooling systems to remove heat from all or a portion of the electronic components.\nAs the trend is to make electronic components more powerful, there is also an increasing push to reduce the size of the electronic components, and the packaging of the electronic components. The smaller components and packaging makes the removal of the unwanted heat more difficult.\nIn some applications, direct impingement thin-film evaporative spray cooling is preferred in order to provide sufficient cooling, whereas in other applications spray cooling is desired to reduce the overall package or housing size even though the required cooling capability is not as high. This creates a situation in which transverse narrow gap evaporative spray cooling is advantageous if it can be done to an acceptable efficiency level.\nNarrow gap evaporative spray cooling will preferably provide or spray the coolant from a transverse side of the surface to be cooled (or the surface from which heat is to be transferred). Proper cooling is preferably achieved if a thin liquid film is maintained over the device or electronic component to be cooled, thereby facilitating evaporation of the coolant as heat is transferred from the electronic component. If there is too little flow or coverage of coolant, the liquid layer covering the electronic component will dry out and cause the component to overheat because vapor forced convection will not typically provide sufficient heat transfer. If the flow of coolant to the component is too great, the device will become flooded and may produce hot spots, insufficient cooling and/or failure, because the vapor created from the evaporation may become trapped between the excessive fluid and the impingement surface of the electronic component. This will normally reduce the cooling efficiency. Vapor generated at the surface of the component which receives too much coolant cannot escape effectively and could result in a boiling heat transfer failure mode generally referred to as burnout.\nEven when the volume flux of coolant is properly matched to the heat flux of the device, the excess fluid sprayed within a cavity must generally be managed by the method described in U.S. Pat. No. 5,220,804 to prevent the overflow from adjacent components from interfering and causing flooding type failure conditions.\nIt is therefore an objective of some embodiments of this invention to provide a narrow gap, thin-film, evaporative spray cooling system for cooling one or more electronic components in the narrow gap.\nIt is also an objective of some embodiments of this invention to provide a narrow gap evaporative spray cooling system which improves the cooling characteristics of the system, especially at the entry end of the cooling channel or conduit, and/or reduces the pressure gradient above the surface from which heat is to be transferred.\nIt is also an objective of some embodiments of this invention to provide a housing system which provides improved re-circulation of the vapor for re-introduction of the vapor into the cooling conduit.\nIt is an objective of some embodiments of this invention to provide a re-circulation system which reduces pooling of the liquid portion of the coolant at or near the exit end of the cooling conduit."}
{"text": "The present invention relates to a plate-like plant carrier for purposes of cultivating layers, i.e. plants in their very early growth, and to have seeds germinate; and particularly the invention relates to a plant growth structure on or as part of plant tables in horticulture, or to be used directly in the ground, or in flower boxes or other open growth systems or in closed systems (glass houses), all for cultivating plants in their early stages.\nIt is known to cultivate plants in boxes, in flats, i.e. in flat boxes made of wood or synthetic material. The flat configuration may also be the result of a small, usually square or rectangularly-shaped and usually plastic flower boxes. What happens to the boxes after planting has obtained is often not deemed a factor, but for reasons of ecology it is. Reference is not made here to the situation of the direct re-use of the boxes after emptying them for planting; of course such boxes can be used over and over again for the same purpose. But presently it is of interest here are that there are situations in which, for some reason, it is not desirable to disturb the seedlings or plants but to leave the container in place, i.e. the plants are just planted in their container. This may readily result in ecological problems.\nThe principle of containing a plant and root system in a small, box-like container carries with it the following problem: a certain amount of water retention is required so that the root system and surrounding soil will not dry out too quickly. On the other hand, if insufficient drainage is provided, then the plant system may begin to rot or be prone to other forms of sickness.\nIt is also known to use very finely ground peat and to press a certain portion into a rod-like configuration. This peat rod is then surrounded by and held in a net made of synthetic material which is configured in this case in that it can deform elastically. Such a rod or strand is then cut into a suitable length and is used for individual cultivating and planting of seeds and small plants. If the pressed strand is watered it will swell as the surrounding synthetic net expands. The finely ground peat, however, has a very small pore volume when compressed, which has an effect on the growth of the plant. In addition, any capillary action is rather poor in compressed peat. Hence, such finally ground peat is a poor agent as far as rewetting is concerned.\nCertain greening flats are known, for example, for use as roof covers or as embankment covers. In the latter case, the task is to have the roots of the plants traverse their container as fast as possible so that the roots physically combine with the underlying soil. In this way one makes sure that the entire arrangement is rapidly and so-to-speak naturally anchored to the ground.\nOn the other hand, roof covers are usually made of mats which contain both plant food and grass seed. These mats are constructed such that the cover on which the mat rests should not receive much moisture, that is, the mat is comprised, for example, of a synthetic layer upon which are deposited layers of nutrients and other plant food to enable the grass to germinate and grow. These mats are constructed for exactly that particular purpose, namely, for providing a green cover on a roof, and these rather highly, specialized mats are usually not useful in a different environment."}
{"text": "The ultimate goal for a NVSRAM memory design is to work like a regular SRAM memory but with a non-volatilability to store the data after power is removed. This is called the Store operation. There are three kinds of Store operations and each time is performed to write SRAM cell's data into each corresponding Flash cell.\nIn a traditional 12T NVSRAM cell, the Write operation is a 2-step operation that uses Erase as a first step to either increase or decrease Flash cell Vt and is followed Program as a second step to conversely decrease and increase Flash cell's Vt to get the final desired Vts of Vt1 (≧2V) and Vt0 (≦−2V) on a single Flash cell. As a result, for traditional 12T NVSRAM cells a Recall operation is based on a method to detect a wide ΔVt (Vt1−Vt0=4V) between one paired flash transistors. But after a long P/E endurance cycle, the distance of the gap of ΔVt of Erase and Program Vt becomes smaller, thus the 12T NVSRAM cell operation becomes critical.\nTherefore, an improved NVSRAM cell design with reduced cell size and proper write operation and recall operation are desired and become objectives of the present invention."}
{"text": "(a) Field of the Invention\nThis invention generally relates to an internal combustion engine that uses rotors or turbines to convert energy released from combustion to shaft power or mechanical output. More particularly, but not by way of limitation, to a rotary motor that includes at least one combustion chamber that delivers the products of combustion in a controlled manner to a turbine or rotor.\n(b) Discussion of Known Art\nThe advantages of a rotary type motor have long been recognized. However, the development of a motor that is able to reliably deliver power harnessed from combustion has seen few successful solutions. The need or absence of such an engine has been particularly acute in the area of smaller turbines where the flow of gasses through the turbine must be carefully controlled in order to prevent damage to the surrounding areas from the exhaust.\nThe use of a separate combustion chamber, which allows the generation of a working fluid is found in the field of steam turbines, where the steam that is used as the working fluid that moves the turbine is generated in a separate boiler that feeds steam to the turbine. This type of system takes advantage of the ease with which steam pressure can be controlled to deliver a desired flow rate of steam to the turbine. Thus, the steam turbine model provides a good model for applications where a working fluid such as steam can be released at a controlled rate. However, this model has proven to be inappropriate for applications where the turbine is to be rotated by the products of combustion of a fuel and oxidant mixture, for example.\nMany known turbine applications that use products of combustion in order to turn the blades of a turbine use a series of turbines that harvest the power released from the rapid expansion from combustion to harvest the needed power to compress the fuel and oxidant mixture and provide power to carry out useful work. Thus, these systems employ a set of axially positioned turbine components that depend on high speed rotation of the components, and very high flow rates of the gases used to power the turbines, in order to provide sustained combustion and rotation of the compressor and power generation components.\nFurthermore, it is well recognized that turbines offer significant advantages over reciprocating engines, particularly in applications where high revolutions per minute (RPMs) are required. However, it has not been practicable to provide a small turbine that can be used in, for example, automotive applications. A significant drawback to the use of turbines in these applications is the high flow rate of gasses generated and discharged during the operation of the turbine. The temperature and flow rate of the exhaust gases delivered from such a turbine create the possibility of causing serious injury to people and property in the vicinity of the exhaust from the vehicle.\nThus, there remains a need for a turbine based engine that can be driven by gases produced from internal combustion, and that permits a highly controlled flow of these gases through the turbine.\nThere remains a need for a turbine base engine that uses at least one turbine to harvest the power released upon combustion, and which allows control of the combustion and delivery of the products of combustion to the turbine.\nIt has been discovered that the problems left unanswered by known art are solve by providing a turbine based engine that includes:\na mixture chamber;\na combustion chamber;\na stationary flow control barrier; and\na turbine. The stationary flow control barrier is located between the combustion chamber and the turbine. It has been discovered that this arrangement allows the system to provide precisely measured and mixed fuel and oxidant mixtures to the mixture chamber where it is burned to produce the expandable gas product that will be expanded through the turbine to turn the turbine.\nAccording to a highly preferred embodiment of the invention the stationary flow control barrier consists of a gas flow control plate that includes ducts that permit gases to flow from the combustion chamber towards the turbine. In a highly preferred embodiment of the invention, these ducts have a tapered or nozzle shaped contour, starting as an aperture of a first dimension near the combustion chamber and progressing to a second, smaller dimension near the turbine. This reduction in the size of the ducts will have the effect of accelerating the flow of the products of combustion as it progresses towards the turbine.\nIt is further contemplated that the fuel and oxidant mixture will be delivered to the combustion chamber by way of valves that control the flow of mixture gases into the combustion chamber. It is contemplated that these valves may control the flow of a fuel and oxidant mixture or the flow of at least one of the components for the mixture and the remaining components for the mixture provided by of an injector or other delivery device.\nAccording to a highly preferred example of the invention, a generally symmetrical arrangement is provided. In this arrangement an expeller device is position along a place of symmetry for the system. Thus, a turbine is positioned on each side of the expeller. The turbines will allow expansion towards the expeller, and thus allow mounting of the expeller and turbines along a single shaft between a pair of mixture chambers, combustion chambers, and flow control barriers.\nIt is further contemplated that these examples may include gas flow control systems that are indexed from the shaft that supports the turbines. It is contemplated that these systems may include cams, markers or other triggering mechanisms that may be used to control a signal that is used to control the delivery of the gasses needed for combustion.\nThus, it should also be understood that while the above and other advantages and results of the present invention will become apparent to those skilled in the art from the following detailed description and accompanying drawings, showing the contemplated novel construction, combinations and elements as herein described, and more particularly defined by the appended claims, it should be clearly understood that changes in the precise embodiments of the herein disclosed invention are meant to be included within the scope of the claims, except insofar as they may be precluded by the prior art."}
{"text": "1. Field of the Invention\nThe invention relates to an air purification system. More specifically, the invention relates to such a system which includes an ozone generator means for producing ozone to react with impurities in the air, and which system also includes a high power corona generating device for removing ozone from the purified air.\nThe invention also relates to a high power corona generating device, and to an air purification system using only the high power corona generating device.\n2. Description of Prior Art\nLow power ozone generators are now being used as air purifiers to treat air in a variety of commercial applications where air quality is unsatisfactory and where special air treatment is desired. Ozone in present purifiers is created when air is drawn through an electrical corona formed over generator plates at high voltage. Such purifiers rely for their effect on adding ozone to the air and allowing natural diffusion to spread the air and ozone through the area to be treated. The concentration of ozone in atmospheres being treated by present purifiers is low and therefore provides sub-optimal results. Some of the disadvantages of the above-described systems are as follows:\nAir throughput into the machine is low so that only a very insignificant amount of air treatment occurs within the machine.\nHealth regulations require that the concentration of ozone in the ambient air of the workplace should not exceed 50-100 parts per billion. Since this low concentration has very limited air purifying capability, this places a very serious constraint on the efficiency and power of present purifiers used in the manner as above-described.\nAir purification systems are also described in U.S. Pat. No. 4,451,435, Holter et al, May 29, 1984, U.S. Pat. No. 3,949,055, Schneider et al, Apr. 6, 1976, and U.S. Pat. No. 4,101,296, Lowther, July 18, 1978.\nThe '435 Patent illustrates, in FIG. 3 thereof, a system including an ozonizer (ozone generator) 11 and beds of sorption masses and catalyst masses 1-9. In accordance with the '435 Patent, the air purification process requires the application of heat and the subsequent application of cooling to remove impurities from the air. Additionally, the apparatus of the '435 Patent does not include any means for decomposing ozone remaining in the air.\nThe '055 Patent teaches the use of an aqueous solution of ozone to remove gaseous or smoke-like substances from the air in industrial areas. The '296 Patent teaches an ozone decomposition system. In neither case could large volumes of air be treated economically.\nA further air purifying system is illustrated in U.S. Pat. No. 4,049,400, Bennett et al, Sept. 20, 1977. The system as taught by Bennett et al includes an ozone generator and an electrostatic filter. This again does not include any means for decomposing ozone remaining in the air.\nIt is also known that certain levels of ozone in the atmosphere can be harmful to people, and that even relatively low levels of ozone can have a corrosive nature. Accordingly, air purifying systems which use ozonators to purify the air but which leave the ozone in the air to be spread throughout the environment of the system could be potentially unsatisfactory.\nOzone generators, per se, are also taught in the following patents which constitutes only a partial list of such teachings. U.S. Pat. No. 4,666,679, Masuda et al, May 19, 1987, U.S. Pat. No. 3,903,426, Lowther, Sept. 2, 1975, U.S. Pat. No. 1,169,825, Hoofnagle, Feb. 1, 1916, U.S. Pat. No. 2,822,327, Hammesfahr et al, Feb. 4, 1958, U.S. Pat. No. 3,607,709, Rice, Sept. 21, 1971, U.S. Pat. No. 3,838,290, Crooks, Sept. 24, 1974, U.S. Pat. No. 4,650,648, Beer et al, Mar. 17, 1987, and U.S. Pat. No. 4,690,803, Hirth, Sept. 1, 1987."}
{"text": "1. Field of the Invention\nThe invention relates to paper machines and in particular to a wet forming section of a paper machine and the operation thereof.\n2. Description of Related Technology\nA paper machine wet forming section having a head box with at least one nozzle for discharging pulp in a jet in a width approximately equal to the width of the machine is known in the art. Such a machine may include a forming cylinder disposed downstream of the head box with respect to the direction of conveyance of the paper pulp through the machine and two continuous-loop forming wires, each wire wrapping about a selected peripheral portion of the forming cylinder, the wires forming a wedge-shaped inlet gap for receiving a jet stream of pulp from the head box nozzle. Such a machine may further include a plurality of guide rolls, one of which being an inlet guide roll around which one of the forming wires is wrapped. It is known to include a displacement device adapted to change the position of the inlet guide roll in order to change the location or length of the peripheral portion of the forming cylinder wrapped by both of the forming wires.\nA paper machine twin-wire web forming system disclosed in Parker et al., U.S. Pat. No. 3,726,758 (Apr. 10, 1973) includes a wedge-shaped inlet gap or nip defined by two continuous-loop forming wires. A pulp jet stream is discharged from a nozzle of a head box into the wedge-shaped gap. FIG. 2 of the Parker et al. patent shows such a device in which the extent that one of the forming wires wraps about a forming cylinder can be altered, as well as the angle the pulp jet is directed to the inlet gap. When the nozzle of the head box is directed appropriately, the pulp jet can be injected either exactly in the middle of the inlet gap, or it can deviate so that the jet stream selectively impinges one or the other wire. As a result, a preliminary dewatering of the pulp occurs on the impinged wire before dewatering is performed between the two wires. The swiveling of the nozzle of the head box serves to continuously change the angle of the pulp jet stream.\nWolf, U.S. Pat. No. 3,944,465 (Mar. 16, 1976) discloses a head box for a paper machine forming section having only a single continuous-loop forming wire or screen. With reference to a side view thereof, the head box swivels in such a way that an opening of a nozzle of the head box moves along a circular arc when the head box is swiveled. The circular arc has a curvature opposite to that of a surface of a breast roll around which the wire is wrapped. The angle of a pulp jet discharged from the nozzle opening can therefore be continuously changed so that the pulp jet stream leaving the nozzle will be discharged either parallel to the wire or impinge the wire at an acute angle. If necessary, the stream is directed upwardly at an acute angle so that the jet is a parabolic trajectory.\nBoth the Parker et al. and Wolf patents disclose ways of influencing dewatering in the first phase of sheet formation. However, such devices have not proved satisfactory as they have been found to also alter the quality of the paper web formed."}
{"text": "1. Field of the Invention\nThe invention relates to the monitoring and analysis of digital information. A method and device are described which relate to signal recognition to enhance identification and monitoring activities.\n2. Description of the Related Art\nMany methods and protocols are known for transmitting data in digital form for multimedia applications (including computer applications delivered over public networks such as the internet or World Wide Web (“WWW”). These methods may include protocols for the compression of data, such that it may more readily and quickly be delivered over limited bandwidth data lines. Among standard protocols for data compression of digital files may be mentioned the MPEG compression standards for audio and video digital compression, promulgated by the Moving Picture Experts Group. Numerous standard reference works and patents discuss such compression and transmission standards for digitized information.\nDigital watermarks help to authenticate the content of digitized multimedia information, and can also discourage piracy. Because piracy is clearly a disincentive to the digital distribution of copyrighted content, establishment of responsibility for copies and derivative copies of such works is invaluable. In considering the various forms of multimedia content, whether “master,” stereo, NTSC video, audio tape or compact disc, tolerance of quality will vary with individuals and affect the underlying commercial and aesthetic value of the content. It is desirable to tie copyrights, ownership rights, purchaser information or some combination of these and related data into the content in such a manner that the content must undergo damage, and therefore reduction of its value, with subsequent, unauthorized distribution, commercial or otherwise. Digital watermarks address many of these concerns. A general discussion of digital watermarking as it has been applied in the art may be found in U.S. Pat. No. 5,687,236 (whose specification is incorporated in whole herein by reference).\nFurther applications of basic digital watermarking functionality have also been developed. Examples of such applications are shown in U.S. Pat. No. 5,889,868 (whose specification is incorporated in whole herein by reference). Such applications have been drawn, for instance, to implementations of digital watermarks that were deemed most suited to particular transmissions, or particular distribution and storage mediums, given the nature of digitally sampled audio, video, and other multimedia works. There have also been developed techniques for adapting watermark application parameters to the individual characteristics of a given digital sample stream, and for implementation of digital watermarks that are feature-based—i.e., a system in which watermark information is not carried in individual samples, but is carried in the relationships between multiple samples, such as in a waveform shape. For instance, natural extensions may be added to digital watermarks that may also separate frequencies (color or audio), channels in 3D while utilizing discreteness in feature-based encoding only known to those with pseudo-random keys (i.e., cryptographic keys) or possibly tools to access such information, which may one day exist on a quantum level.\nA matter of general weakness in digital watermark technology relates directly to the manner of implementation of the watermark. Many approaches to digital watermarking leave detection and decode control with the implementing party of the digital watermark, not the creator of the work to be protected. This weakness removes proper economic incentives for improvement of the technology. One specific form of exploitation mostly regards efforts to obscure subsequent watermark detection. Others regard successful over encoding using the same watermarking process at a subsequent time. Yet another way to perform secure digital watermark implementation is through “key-based” approaches."}
{"text": "Various means have been developed for use in analyzing characteristics of lubricating oils, including engine oils. In particular, methods and apparatus for testing the condition of oil and the sludge content of oil include chromatography and chemical analysis.\nOther methods and apparatus for assessing the quality of used oil include placing a measured amount of oil upon an absorbent material, heating the sample, and awaiting dispersion of the sample. The amount of undispersed sludge may then be measured and rated quantitatively. These methods and apparatus, however, require significant controlled conditions, including measurement of the oil sample volume, and the use of a template to measure and rate the quantity of undispersed sludge in the sample. Additionally, these methods include heating of the sample, and awaiting dispersal of the sample. A need exists for a simple and rapid method of analyzing an oil sample on a qualitative basis."}
{"text": "Video games and video game systems have become extremely popular. Video game devices or controllers typically use visual and auditory cues to provide feedback to a user. In some interface devices, kinesthetic feedback (e.g., active and resistive force feedback) and/or tactile feedback (e.g., vibration, texture, temperature variation, and the like) may be provided to the user. In general, such feedback is collectively known as “haptic feedback” or “haptic effects.” Haptic feedback provides cues that enhance and simplify a user's interaction with a video game controller, or other electronic device. For example, haptic effects may provide cues to users of video game controllers or other electronic devices to alert the user to specific events, or provide realistic feedback to create greater sensory immersion within a simulated or virtual environment.\nOther devices in which a user interacts with a user input element to cause an action also may benefit from haptic feedback or haptic effects. For example, such devices may include medical devices, automotive controls, remote controls, and other similar devices."}
{"text": "1. Field of the Invention\nThe invention relates to a display device constituted including a transistor and a driving method of the display device. In particular, the invention relates to a semiconductor device including a pixel constituted including a thin film transistor (hereinafter also called a transistor).\n2. Description of the Related Art\nAn active matrix display, which is constituted by the combination of an electroluminescence element (also called an organic light emitting diode (OLED) and an EL element or a light emitting element in this specification) and a transistor, has been attracting attentions and actively researched and developed both domestically and internationally as a thin and lightweight display. This display which is also called an organic EL display (OELD) is extensively researched and developed in a practical use stage aiming for a small 2-inch display to a large display of 40-inch or larger.\nLuminance of an EL element and a current value flowing therethrough are theoretically in a linear relationship. Therefore, for an organic EL display which employs an EL element as a display medium, a method to express a gray scale by controlling a current value supplied to the EL element is known. Moreover, as a method to control a current value supplied to the EL element, a voltage input driving method and a current input driving method are known.\nIn the voltage input driving method, a current value supplied to a driving transistor (hereinafter also called a driving transistor) and an EL element is controlled by a gate-source voltage obtained by inputting a voltage signal to a gate of a driving transistor so as to be held therein, which is connected in series to the EL element. In the current input driving method, a current value supplied to a driving transistor and an EL element is controlled by a gate-source voltage of a driving transistor obtained by supplying a current signal to the driving transistor (for example, refer to Patent Document 1).\nHowever, in a conventional current input driving method, a slight amount of current is required to be supplied from a source signal line to express a low gray scale. As time to charge parasitic capacitance of a source signal line or the like is required to input a slight amount of current as a video signal to a pixel, there is a problem in that long writing time is required.\nFurther, as another example of a current input driving method, such a pixel is known, in which by holding Vgs inputted as a current to a driving TFT and a threshold voltage thereof in two capacitors and capacitively coupling them, a current supplied to an EL element can be smaller than an actual video signal while compensating the threshold voltage (for example, refer to Patent Document 2).\nHowever, even such a pixel configuration requires a period T1 to obtain a threshold voltage and a period T2 to write a video signal. As an area of one pixel is limited, the capacitance of two capacitors are also limited. Therefore, there is a problem in that there is not enough writing time for writing a slight amount of current as a video signal, and in a large panel in particular, a writing period per pixel becomes shorter as compared to a small panel.\n[Patent Document 1]\n    International Publication No. 9848403[Patent Document 1]    Japanese Patent Laid-Open No. 2004-310006"}
{"text": "The Common Object Request Broker Architecture (CORBA) middleware platform has been a leading middleware platform in recent years. As is known in the art, CORBA is a standard defined by the Object Management Group (OMG) that enables software components written in multiple computer languages and running on multiple computers to work together as a single application or set of services. CORBA uses an interface definition language (IDL) to specify interfaces that objects will present to the outside world, and specifies a mapping from the IDL to a specific implementation language. CORBA uses an Object Request Broker (ORB) to send requests from objects executing on one system to objects executing on another system. The ORB allows objects to interact in a heterogeneous, distributed environment, independent of the computer platforms on which the various objects reside and the languages used to implement them. CORBA is specified and further explained in the CORBA Specification, version 3.1 (January 2008), including Part 1 (CORBA Interface) and Part 2 (Interoperability), available from the Object Management Group (OMG) at 109 Highland Ave, Needham, Mass. 02494. The entire CORBA Specification, including at least Parts 1 and 2, is hereby incorporated by reference in its entirety.\nCORBA communication typically occurs over an Ethernet or other network, between a Client and Server. FIG. 1 is an illustrative prior art block diagram 10 showing client 12 to server 14 communication, in accordance with CORBA, in an exemplary environment. FIG. 1 illustrates what happens when a client application 16 invokes an operation on an object/servant 24 in a server process 14. To implement an interface, CORBA IDL is compiled into the source code language with which the client 12 or server 14 is implemented. On the client side, this code is called a stub. On the server-side, this IDL code is called a skeleton. Typically, client-side application code 16 invokes a local proxy object 18 (e.g., via a proxy class generated by an IDL compiler). The proxy 18 gets information about the request (e.g., in and inout parameters, operation name) into a binary buffer, which is then passed into the ORB 20A library. The ORB 20A library sends a request message across the network to the server process 14. The ORB 20A waits for a reply message from the server process 14. The ORB 20A returns the reply buffer back to the proxy object 18, which unmarshals inout and out parameters and the return value (or a raised exception), and returns these to the client application code 16.\nAt the server 14 side, the ORB 20B runs a thread in an event loop that waits for incoming requests. When the request arrives from the client 12 the ORB 20B reads the request's binary buffer and passes this to some code that unmarshals the parameters and dispatches the request to the target servant 24. The code that performs the unmarshalling and dispatching is spread over two components, the Portable Object Adapter (POA) (shown in FIG. 1) and the skeleton code that is generated by the IDL compiler. When the operation in the servant 24 returns, the skeleton code marshals the inout and out parameters (or a raised exception) into a binary buffer and this is returned via the POA 22 to the ORB 20B, which transmits the reply message across the network to the client process 12.\nThe main protocol for ORB communication as shown in FIG. 1 is the standardized General Inter-ORB Protocol (GIOP), which has been widely deployed for transport in the TCP/IP environment. GIOP also is described further in the aforementioned CORBA Specification, and as of this writing is a version 1.3. GIOP over TCP/IP is known as Internet Inter-ORB Protocol (IIOP). GIOP is a client-server protocol and defines the messages and format that are passed over the ORB between the client and the server object. The data placed in the GIOP follows CORBA Common Data Representation (CDR) syntax for placing and copying the data into an octet stream. CORBA is mainly used on General Purpose Processors (GPPs) using TCP/IP and/or OS Inter-Processor Communications (e.g., shared memory).\nThere is increasing need to apply technology such as the CORBA GIOP ORB in different types of environments, but some environments, especially embedded environments, require more efficient and/or compact messaging than is provided via the GIOP message and GIOP header formats of FIG. 2. In addition, standard GIOP/IIOP interoperability protocols can be less than optimal for applications having strict requirements for latency, overhead and message sizes. Furthermore, because the CORBA GIOP was originally developed for use in general purpose distributed computing environments, optimization of the GIOP may be required for the best performance in distributed embedded systems, which can be more complex, especially because of the many interfaces with different types of control devices and input/output (I/O) devices. Optimized interoperability protocols are thus becoming of greater importance."}
{"text": "The present disclosure relates generally to plastic (sheet) laminates, and more specifically to a plastic and glass laminate with a filling therebetween.\nGlass laminated products have contributed to society for almost a century. Beyond the well known, every day automotive safety glass used in windshields, glass laminates are used in most forms of the transportation industry. They are utilized as windows for trains, airplanes, ships, and nearly every other mode of transportation. Safety glass is characterized by high impact and penetration resistance and does not scatter glass shards and debris when shattered. Glass laminates find widespread use in architectural applications, as well.\nA glass laminate typically consists of a sandwich of two glass sheets or panels bonded together with an interlayer of a polymeric film or sheet which is placed between the two glass sheets. One or both of the glass sheets may be replaced with optically clear rigid polymeric sheets such as, for example, sheets of polycarbonate materials. Glass laminates have further evolved to include multiple layers of glass and/or polymeric sheets bonded together with interlayers of polymeric films or sheets.\nAlthough the glass-plastic laminate has the advantage of a reduced weight compared to a glass-glass laminate, the glass-plastic laminate can be stressed during thermal cycles and as a result has issues such as curvature and/of cracking over time. Hence there is a continual need for glass-plastic laminates with reduced stress built-up during thermal cycles and enhanced resistance to cracking."}
{"text": "The present invention concerns an electronic mixer for generating a mixed signal by mixing a local oscillator signal with a useful signal comprising at least one field effect transistor which has at least one gate, at least one source and at least one drain, a useful signal input for feeding in the useful signal with a useful frequency, a local oscillator signal input for feeding in the local oscillator signal, which is so designed that in operation of the mixer it receives a local oscillator signal whose local oscillator frequency is equal to an integral fraction of the useful frequency which is reduced or increased by a mixing frequency, and a signal output at which the mixed signal is present in operation of the mixer.\nThe present invention further concerns a method of generating a mixed signal by mixing a local oscillator signal with a useful signal in at least one field effect transistor having at least one gate, at least one source and at least one drain, comprising the step of generating the local oscillator signal with a local oscillator frequency, wherein the local oscillator frequency is equal to an integral fraction of the useful frequency which is reduced or increased by the mixing frequency.\nFor efficiently acquiring electromagnetic signals beat receivers, in particular heterodyne receivers, have long been used, in which the electromagnetic signal to be acquired, which is referred to hereinafter as the useful signal, is mixed with a signal which is generated locally, that is to say at or in the mixer itself, which hereinafter is referred to as the local oscillator signal. The mixed signal then has a mixing frequency which is equal to the differential frequency between the frequency of the useful signal, this is referred to hereinafter as the useful frequency, and the frequency of the local oscillator signal, this is referred to hereinafter as the local oscillator frequency. The amplitude of the mixed signal is a measurement in respect of the amplitude of the useful signal. The signal identified as the mixed signal with the mixing frequency in the present application as the output of the mixer is frequently also referred to in the literature as the intermediate frequency signal with the intermediate frequency.\nThe terahertz frequency range or submillimeter wavelength range which is roughly defined by between 100 gigahertz (GHz) and 10 terahertz (THz) is one of the last ‘dark’ areas of the electromagnetic spectrum.\nTechnically usable detectors are not commercially available in that frequency range or are commercially available only at low frequencies. In particular the mixing efficiency of field effect transistors which are frequently used for mixing a useful signal with a local oscillator signal drops off severely towards higher useful frequencies.\nIn that respect however the efficiency of the mixers also depends on the available power of the local oscillator signal. Powerful local oscillator sources are not available at high terahertz frequencies, or are so available only at a high level of technical and financial complication and expenditure. Therefore it is known from the state of the art to use a local oscillator which generates a local oscillator signal whose local oscillator frequency is equal to an integral fraction of the useful frequency increased or reduced by the mixing frequency. Such a mixing method is referred to as subharmonic mixing.\nIn subharmonic mixing, the non-linear properties of the mixer component, thus also in the case of the field effect transistor, are put to use. A mixed product is generated with the difference of an integral multiple of the local oscillator frequency and the useful frequency.\nIt will be noted however that the efficiency of subharmonic mixing is markedly lower than the mixing efficiency upon coupling the fundamentals involving the same power into the mixer element, in particular the field effect transistor."}
{"text": "1. Field of the Invention\nThe present invention relates to a shaver head for a wet shaver. The shaver head includes a razor blade and is arranged at the front end of a handle. The shaver head also includes a guide strip parallel to and in front of the cutting edge of the razor blade.\n2. Description of Related Art\nConventional shaver heads for a wet shaver are arranged at the front end of a handle and include a single or double razor blade, which is covered by a cover cap. A guide strip is arranged parallel to and in front of the cutting edges of the razor blades.\nIn such a shaver head, the shaving geometry, and thus the shaving angle, are defined by the front guide strip, the rear cover cap and the cutting edges of the razor blades. However, these parts are fixed to one another, resulting in a completely predetermined shaving geometry and shaving angle, upon which the shaving properties depend. It is desirable, however, for the user to be able to vary the shaving geometry, and thus the shaving angle in accordance with his requirements."}
{"text": "This invention relates to a punch for punching markings into rolled material such as bars, ingots, billots or slabs, relative to which the rolled material is guided, including a toolholder movable by means of a controlled piston-cylinder arrangement, the punch types being settable automatically.\nSuch apparatus are known (see for instance U.S. Pat. No. 3,306,186). They are provided with a plurality of punch plates which are rotatably mounted on a common shaft guided in a toolholder. The toolholder is firmly arranged on a slide which is movable by means of a controlled reciprocating cylinder. In such apparatus all punch markings can be produced only simultaneously, i.e. in one operation, in the event the marking to be provided on the rolled material consists of a plurality of letters, figures or other symbols. This brings about the essential disadvantage that in the event of a non-planar end face or cutting surface of the rolled material or in the event of an inclined position of the rolled material in relationship to the punch plates uniform punch markings readily to be discriminated are not insured. An apparatus is already known wherein the slide carrying the toolholder with its slideway as well as the reciprocating cylinder associated with the slide are firmly arranged at the lower side of a support plate horizontally pivotable in counteraction to the lateral pressure of helical springs and wherein at the slideways at the end face two spaced, horizontally arranged rod-shaped feelers projecting beyond the end face of the slideways are positioned which upon engaging the rolled material set the punch plates via the support plate parallel to the surface of the rolled material to be provided with the marking. A proper punch marking is achieved with this apparatus only, however, in the event that the end face of the rolled material is substantially planar. When the end face is non-planar, the punch plates are only able to engage in varying depths in spite of a setting parallel to the end face of the rolled material. A further disadvantage inherent to the conventional apparatus is that the marking regarding the number of individual punch types is dependent on the number of punch plates with which the apparatus is equipped, i.e. for punching for instance a five-digit number five punch plates are required. Furthermore, in the prior art apparatus the spacing of the individual punch markings relative to one another is not variable."}
{"text": "In a time-of-flight (TOF) range sensor configured to acquire a range image, by using TOF scheme, a potential just under a gate electrode of a MOS structure is controlled in a vertical direction (depth direction) of the MOS structure. For example, a CMOS distance-measuring element and a TOF image sensor using the CMOS distance-measuring element are disclosed in patent literature (PTL) 1. The CMOS distance-measuring element has a structure such that an n-type charge-generation buried region buried in a p-type semiconductor layer, a charge-transfer buried region, a charge-read-out buried region, an insulating film covering the charge-generation buried region and the charge-transfer buried region, a transfer gate electrode provided on the insulating film to transfer signal charges to the charge-transfer buried region, a read-out gate electrode provided on the insulating film to transfer the signal charges to the charge-read-out buried region. When pulse lights are irradiated to the charge-generation buried region in the CMOS distance-measuring element recited in PTL 1, light signals are converted to signal charges in a semiconductor layer just under the charge-generation buried region, and a distance to an object is measured from a distribution ratio of charges accumulated in the charge-transfer buried region.\nThe CMOS distance-measuring element or the TOF image sensor using the CMOS distance-measuring element has a problem of generation of noise or dark current caused by interface defects, interface states, or the like just under the transfer gate electrode. In addition, in the case of using the transfer gate electrode disclosed in PTL 1, a potential gradient over a long distance is difficult to control, and a substantially uniform electric field over a long distance of a charge transport path is practically difficult to maintain. For this reason, in the charge-modulation elements such as distance-measuring elements having long charge transport paths, carriers are stopped in the middle of the charge transport paths, and thus, there are issues such that expected performances of the charge-modulation element are difficult to achieve."}
{"text": "In recent years, a portable device which can display a video with high definition such as a smartphone and a tablet has been widely used. Accordingly, development in the MHL which is a communication interface standard to transmit a video at high speed for the portable device has been proceeded (for example, refer to Patent Document 1).\nAs a communication interface standard for realizing uncompressed digital video transmission, the high definition multimedia interface (HDMI®) (registered trademark) is exemplified. Whereas, main characteristics of the MHL is to minimize a mounting area as a minimum pin configuration necessary for video transmission and to assist power supply.\nThe MHL devices are classified into three categories, i.e., a source device for transmitting a video signal, a sink device for receiving and displaying the video signal, and a dongle device for converting the video signal in the MHL format into the other video signal. Then, an MHL cable which satisfies the MHL standard is used to connect the MHL devices and to transmit signals between the MHL devices. The source device includes a personal computer, a smartphone, a tablet terminal, a game machine, and a digital camera. Also, the sink device includes a display device such as a digital TV. A single MHL cable connects the source device to the sink device so that a video with high definition can be transmitted and power can be supplied (charge the source device).\nIn a communication system according to the MHL standard, basically, regular power supply is started after a link is established between the source device and the sink device. Therefore, when the MHL cable used to connect between the source device and the sink device is an active cable and when the source device (Direct Attached Device) is directly connected to the sink device without using the MHL cable, there is a problem in that a failure such that start-up is not available due to short power supply at the time of start-up is caused or that it is necessary to have an external power supply."}
{"text": "1. Field of the Invention\nThe present invention relates to an automatic answering method and apparatus for supporting a question reply process of replying to a question document of a text format.\n2. Description of the Related Art\nWith recent widespread of computerization, questions to companies or the like are often made by form inputs at home pages or e-mails. If every question is to be answered manually on the company side, many operators are required and the cost increases. A novice operator can not answer some questions or it takes a long time for the novice operator to answer a question. In order to solve this problem, a question-answering system has been introduced recently. With this system, a question document is input and its content is analyzed to select a reply example candidate from reply examples and question-reply examples prepared for each question content and to present the selected reply example candidate.\nMost of such question-answering systems assume, however, that one document contains only one consultation content. Therefore, if a plurality of question contents are written in one document, the systems cannot analyze each question content, resulting in a low reply precision.\nAnother technique is disclosed in JP-A-2002-132661. This technique discloses means for dividing one document containing a plurality of question contents, into each question content. The divided question content is analyzed to select a reply example candidate. A reply precision representative of a likelihood or degree of each reply example candidate for the question content is calculated. If the reply precision has a predetermined value or higher, an answer is formed from the reply example candidate, whereas if the reply precision is lower than the predetermined value, an instruction is given to compose a new answer.\nThe conventional technique disclosed in JP-A-2002-132661 describes that the means for dividing a document into each question content performs a division process by using “number”, “alphabet”, “.”, an indent, a conjunction such as “or”, and the like. However, if a document is divided into each question content by using “number”, an indent and the like as a separator, there occurs the problem that one question content is divided into a plurality of sentences. Conversely, there arises the problem that if the range of a question content is broad, example candidates for a plurality of question contents cannot be selected.\nAccording to conventional techniques, since a question document is divided basing upon only the information about the contents of the question document, the divided range may not be covered by each reply example candidate. Namely, it is necessary to divide a question document so as to be covered by a prepared reply example candidate, and not to divide it by referring only to the question document content.\nSince a question document divided basing upon conventional techniques may be a document irrelevant to the question document content, the reply example candidate generation process is adversely affected so that the reply example candidate generation precision lowers. It also takes a time for a reply composition operator to find a proper document to be read.\nAccording to conventional techniques, a reply precision representative of the likelihood value of a reply example candidate is calculated, and if the reply precision is a predetermined value or higher, a reply is generated from the reply example candidate to automatically answer (automatically return) the question. If the reply precision is lower than the predetermined value, an instruction is given to compose a new answer. However, if there are a large number of types of replies or if a similar question requires a different answer, the reply precision lowers so that the number of samples exceeding a predetermined threshold reduces. Therefore, the number of samples capable of being used for the automatic reply reduces, and the number of cases requiring to generate new answers increases. There arises the problem of a low operator work efficiency or an automatic reply using an erroneous reply example candidate."}
{"text": "In the past a number of reversing traverse mechanisms have been proposed and constructed. In a common friction type device, one or more roller elements were swivel-mounted on a base member which moved along a rotating shaft at a given speed, depending on the rate of shaft rotation and the angle at which the roller means was disposed with respect to the shaft. Frequently, several rollers were used to engage the shaft in order to minimize slippage and increase efficiency. In case where more than one roller was employed, they were usually connected by means of gears or linkages to enable all of them to be swiveled simultaneously in response to swiveling of any given one. A reversal in the direction of travel of the mechanism was effected by engagement of a swivel arm connected to one roller with a fixed abutment located at a predetermined point along the shaft. However, there were distinct disadvantages to this construction, in that as the swivel arm was turned, the relative angle between the rollers and the shaft slowly decreased which resulted in a slow-down of the traversing movement of the mechanism immediately before reversal. Under some circumstances there was a distinct possibility that as the rollers were swiveled they could, at some point during a reversal, assume a \"dead\" position wherein their axes of rotation would be parallel to the shaft axis. Under such a condition the traversing movement would cease and the rollers would have to be manually swiveled to initiate further traversing movement."}
{"text": "The present invention concerns new solid solutions which are useful as inorganic color pigments. More particularly, the present invention concerns new solid solutions having a corundum-hematite crystalline structure which are useful as inorganic color pigments, some of which exhibit low Y CIE tri-stimulus values and high reflectivity in the near infrared portion of the electromagnetic spectrum.\nChromium green-black hematite (basic chemical formula: Cr2O3) is an inorganic color pigment, C.l. Pigment Green 17, having a corundum-hematite crystalline structure. It is commonly used to impart a green color to ceramics, paints, polymers, and other materials. The DCMA Classification and Chemical Description of the Complex Inorganic Color Pigments, Third Addition (1991), published by the Dry Color Manufacturer\"\"s Association, states that its composition may include any one or a combination of the modifiers Al2O3 (alumina), Fe2O3(iron oxide), or Mn2O3(manganese oxide).\nIron brown hematite (basic chemical formula: Fe2O3), is an inorganic pigment, C.l. Pigment Red 101 and 102, having a dark brown to black color and a corundum-hematite crystalline structure. See DCMA Classification and Chemical Description of the Complex Inorganic Color Pigments, Third Addition (1991). It is commonly used to impart dark brown to black color to ceramics, paints, plastics, and other material. Its composition may include any one or a combination of the modifiers Cr2O3(chrome oxide), Fe2O3(iron oxide), Mn2O3(manganese oxide), or NiO (nickel oxide).\nChromium green-black hematite is one of the principle pigments used in the manufacture of green shade military camouflage paint and netting. In such applications, chromium green-black hematite is combined with cobalt bearing mixed metal oxides, such as cobalt containing spinel pigment V12600 available from Ferro Corporation of Cleveland, Ohio. This combination of pigments is effective in simulating the reflectivity of chlorophyl in the visible portion of the electromagnetic spectrum, being that portion of the spectrum which is viewable by the naked eye with wavelengths ranging from approximately 0.40 xcexcm to 0.70 xcexcm.\nChlorophyl, which is an organic pigment, generally exhibits a relatively uniform high degree of reflectivity in the near infrared, being that portion of the electromagnetic spectrum with wavelengths ranging from approximately 0.7 xcexcm, to 2.5 xcexcm. Cobalt, however, exhibits a strong absorption band (i.e., low reflectivity) in a portion of the near infrared with wavelengths ranging from approximately 1.2 xcexcm to 1.6 xcexcm. In recent years, advancements in imaging technology have made it possible to contrast known military green shade camouflage painted or covered objects from the background foliage in that portion of the near infrared. A substitute military green shade camouflage pigment which contains no cobalt and which closely matches the reflectivity of chlorophyl in the visible and near infrared is therefore highly desired.\nIn order to satisfy military specifications, a substitute green shade camouflage pigment would have to exhibit a dark drab green appearance in the visible portion of the spectrum and would also have to simulate the reflectance curve for chlorophyl in the near infrared. Generally speaking, known inorganic pigments which exhibit a low degree of reflectivity in the visible portion of the light spectrum (i.e., dark drab colored pigments) also tend to exhibit a correspondingly low degree of reflectivity (i.e., high absorption) in other portions of the light spectrum, including the near infrared. A chromium green-black hematite pigment manufactured by Bayer Corporation of Germany, product number AC 5303, was observed to exhibit a higher near infrared reflectance than other chromium green-black hematite sources (this pigment, however, does not have the desired dark drab appearance in the visible spectrum required for military green shade camouflage paint applications). It was found by chemical analysis that this pigment contained both alumina and titania (basic chemical formula: TiO2) as minor additives. A search failed to disclose any references teaching the use of alumina and titania to improve the near infrared reflectance of chromium green-black hematite pigments.\nThe present invention provides new solid solutions having a corundum-hematite crystalline structure which are useful as inorganic color pigments. Solid solutions according to the present invention are comprised of a host component having a corundum-hematite crystalline structure which contains as a guest component one or more elements from the group consisting of aluminum, antimony, bismuth, boron, chrome, cobalt, gallium, indium, iron, lanthanum, lithium, magnesium, manganese, molybdenum, neodymium, nickel, niobium, silicon, tin, titanium, vanadium, and zinc. Solid solutions according to the present invention are formed by thoroughly mixing compounds, usually metal oxides or precursors thereof, which contain the host and guest components and then calcining the compounds to form the solid solutions having the corundum-hematite crystalline structure.\nSome of the new solid solutions according to the present invention, such as for example chrome oxide as a host component containing the elements iron, boron, and titanium as guest components, exhibit dark drab colors in the visible and high reflectivity in the near infrared portions of the electromagnetic spectrum. One of the primary uses for new solid solutions having these properties would be as inorganic color pigments in military camouflage paint or netting applications, which would permit the radiation signature of a painted or covered object to be tailored to match the reflectance curve of the background foliage in the visible and near infrared portions of the electromagnetic spectrum. Because many of these new solid solutions exhibit relatively high near infrared reflectance, they would also be suitable for use in the general paint and polymer markets, most specifically for architectural applications where increased near infrared reflectance would result in lower heat build-up and thus lower energy costs.\nThe foregoing and other features of the invention are hereinafter more fully described and particularly pointed out in the claims, the following description setting forth in detail certain illustrative embodiments of the invention, these being indicative, however, of but a few of the various ways in which the principles of the present invention may be employed."}
{"text": "Communication systems employ coding to ensure reliable communication across noisy communication channels. These communication channels exhibit a fixed capacity that can be expressed in terms of bits per symbol at certain signal to noise ratio (SNR), defining a theoretical upper limit (known as the Shannon limit). As a result, coding design has aimed to achieve rates approaching this Shannon limit. One such class of codes that approach the Shannon limit is Low Density Parity Check (LDPC) codes.\nTraditionally, LDPC codes have not been widely deployed because of a number of drawbacks. One drawback is that the LDPC encoding technique is highly complex. Encoding an LDPC code using its generator matrix would require storing a very large, non-sparse matrix. Additionally, LDPC codes require large blocks to be effective; consequently, even though parity check matrices of LDPC codes are sparse, storing these matrices is problematic.\nFrom an implementation perspective, a number of challenges are confronted. For example, storage is an important reason why LDPC codes have not become widespread in practice. One of the major issues in LDPC decoding is the organization of memory. Recognizing the fact that the larger the memory size, the lower is the cost of per bit, there is motivation to the investigate of code architectures that allow efficient lumped memory structure for the large amount of edge values in a LDPC decoder.\nAlso, a key challenge in LDPC code implementation has been how to achieve the connection network between several processing engines (nodes) in the decoder. Further, the computational load in the decoding process, specifically the check node operations, poses a problem.\nIt is recognized that, in general, a code corresponding to set number of parallel engines cannot be reconfigured to accommodate different numbers of parallel engines such that an efficient memory architecture is maintained. Because different applications require different decoding speeds, this inflexibility presents a serious drawback. Also, this constraint hinders the ability to adopt and exploit advancements in semiconductor capabilities. For instance, as processing power increases, the number of parallel engines needed for a given decoding speed should reduce; however, the design of a fixed number of engines does not permit a straightforward reduction in the number of processors utilized.\nTherefore, there is a need for a LDPC communication system that employs efficient decoding processes. There is also a need to minimize storage requirements for implementing LDPC coding. There is a further need for a decoding scheme that can readily accommodate future technological advancements."}
{"text": "This invention relates to a novel acrylic two-part type adhesive suitable for structural uses having excellent storage stabilities and adhesion performances.\nHitherto as acrylic adhesives have been known anaerobic adhesives, cyanoacrylate adhesives and the like. They are characterized by exhibiting adhesion in a short time and exhibiting a strong tensile strength although they are of one-pack type, while they are very weak against some strengths applied such as tearing and impact. Thus, their uses are restricted to the fixing of fitted portions or the like in the case of the anaerobic adhesives and to temporary fixing or the like in the case of the cyanoacrylate adhesives.\nOn the other hand, acrylic two-part adhesives wherein an elastomer such as acrylic rubber or epichlorohydrin rubber was dissolved in an acrylic monomer such as methyl methacrylate have been known (sometimes referred to as the first-generation acrylic adhesives). The adhesives of this type, however, do not undergo in a curing step a chemical reaction between the monomer and the elastomer, wherein the fragility caused by the acrylic monomer alone has been simply improved by the presence of the elastomer. Moreover, they exhibit good working performance but are considerably inferrior in adhesion performances, in comparison with the epoxy adhesives which are now used for structural uses in the greatest quantities.\nRecently have been developed acrylic two-part adhesives called \"the second-generation acrylic adhesives\". In these adhesives, a chlorosulfonated polyethylene is often used as the elastomer component. The adhesives are characterized by such a mechanism that radicals are produced on the side chains of the elastomer in the course of curing and polymerizable monomers are graft-polymerized thereto. Thus, the adhesives have excellent adhesion properties in comprison with the first-generation adhesives which do not involve such a graft polymerization. More specifically, the adhesives have the tensile shear strength and impact strength equivalent or even superior to those of the epoxy adhesives, and also exhibit far better values with respect to the tear strength, fatigue strength and oily surface adhesion properties than the epoxy adhesives.\nAs to the working performances, the second-generation adhesives (as well as the first-generation adhesives) are advantageous in that they do not require troublesome operations such as weighing and mixing which are required in the case of the epoxy adhesives. In this respect, the acrylic adhesives of this type may be said to be an excellent adhesive. However, when they are used to bond metal articles, the adhesive containing chlorosulfonated polyethylene as the elastomer is dechlorinated after long term storage or upon heating at a high temperature and the resulting chlorine often corrodes the bonded metal surfaces to lower the adhesion strength. Thus, a small amount of an absorbent for the decomposed chlorine such as epoxy resins has been incorporated to stabilize the adhesive, but a satisfactory effect can not be exhibited. Moreover, a second-generation adhesive wherein a diene-elastomer was used has been developed, but this adhesive is also not satisfactory with respect to its adhesion performances and storage stability.\nThe present inventors have been engaged in the researches on curable compounds having ionic bonds, and invented one-pack anaerobic adhesives having excellent properties. Such adhesives include, for example, a rapidly curing anaerobic adhesive comprising an anaerobic mixture base of a polyvalent metal salt of an acid ##STR1## and an acrylate or methacrylate ester and a small amount of an organic acid adduct of an amine and organic peroxide (cf. Japanese Patent Publication No. 47492/1977); and a one-pack anaerobic adhesive, which is rapidly curable and can be strongly bonded to a metal other than iron and copper without a primer, comprising the above-mention anaerobic mixture base and o-benzsulfimide, tetrahydroquinoline and an organic peroxide (cf. Japanese Patent Publication No. 477766/1977).\nAfter intensive researches, the present inventors have accomplished a novel adhesive of the second-generation acrylic type which is suitable for structural uses and excellent in storage stability and adhesion performances."}
{"text": "A Group 5 metal oxide film such as niobium oxide and tantalum oxide is attracting attention as a film having a high refractive index and exhibiting high optical transparency, i.e., as a material for a device such as semiconductor device and optical device. The technique for manufacturing a Group 5 metal oxide film includes, in rough classification, two processes, i.e., a dry process and a wet process. The dry process includes a sputtering method, an ion plating method, an atomic layer deposition method (ALD method), a chemical vapor deposition method (CVD method), etc. The wet process includes a sol-gel method, an organic metal deposition method (Metal Organic Deposition; MOD method), etc. While the dry process requires a special production facility such as large vacuum apparatus, the wet process has a cost benefit in that the process can be performed only with a simple production facility. In the case of manufacturing a film by a wet process, the kind of the film-forming material and the film formation temperature greatly affect the quality of the film obtained. In Non-Patent Documents 1 to 12, a multinuclear niobium complex and a multinuclear tantalum complex, each having an oxo ligand and an alkoxy ligand, are described. The complexes described in those documents are different from the Group 5 metal oxo-alkoxo complex of the present invention in the number of central metals or the number of ligands. In addition, those documents are absolutely silent as to using the complex described in the documents as a film-forming material."}
{"text": "1. Field\nAn exemplary embodiment relates to a display device, and more particularly, to an organic light emitting device.\n2. Description of the Related Art\nThe importance of flat panel displays has recently increased with the growth of multimedia. Various types of flat panel displays such as liquid crystal displays (LCDs), plasma display panels (PDPs), field emission displays (FEDs), and organic light emitting devices have been put to practical use.\nIn particular, an organic light emitting device may have a high response speed (of 1 ms or less), a low power consumption, and a self-emitting structure. The organic light emitting device may also not have viewing problems. As such, the organic light emitting device has been considered as a next generation display device.\nThe organic light emitting device and an active matrix type organic light emitting device depending on a driving manner. In the passive matrix type organic light emitting device, an anode electrode and a cathode electrode cross each other at a right angle, and signal lines are selected to thereby drive the organic light emitting device. In the active matrix type organic light emitting device, a thin film transistor is connected to each pixel electrode, and the organic light emitting device is driven depending on a voltage maintained by capacitance of a capacitor connected to a gate electrode of the thin film transistor.\nThe organic light emitting device may include a light emitting diode including a first electrode, a light emitting layer, and a second electrode. More specifically, the light emitting diode includes the first electrode supplying holes to each subpixel, the second electrode supplying electrons to each subpixel, and the light emitting layer interposed between the first electrode and the second electrode. The light emitting layer forms excitons by combining the holes received from the first electrode and the electrons received from the second electrode to thereby emit light.\nThe light emitting diode is formed by evaporating a metal having a high work function such as indium-tin-oxide (ITO) and indium-zinc-oxide (IZO) and then patterning the evaporated metal to form the first electrode in each subpixel, and forming a third insulating film that covers a portion of the first electrode and defines each subpixel.\nThe light emitting layer is formed on an exposed area of the first electrode exposed by the third insulating film. The light emitting layer may include an organic material or an inorganic material. In case that the light emitting layer may include an organic material, the light emitting layer formed of the organic material may be formed using a thermal evaporation method in which the light emitting layer is evaporated upward by applying heat to an evaporation source.\nIn the thermal evaporation method, in case that the substrate is close to the evaporation source, an evaporation thickness is uniform. On the contrary, in case that the substrate is far away from the evaporation source, a scattering angle of the evaporation source in an outermost area of the substrate is small. Therefore, an evaporation shadow phenomenon, in which a predetermined area of the light emitting layer is not evaporated, occurs. Accordingly, the evaporation shadow phenomenon reduces the reliability of the organic light emitting device and generates dark spots on an image."}
{"text": "1. Technical Field of the Invention\nThis invention relates to software fault management systems and, more particularly, to a method of correlating multiple network alarms in a large communications network.\n2. Description of Related Art\nIn communications networks, a single network fault may generate a large number of alarms over space and time. In large, complex networks, simultaneous network faults may occur, causing the network operator to be flooded with a high volume of alarms. The high volume of alarms greatly inhibits the ability to identify and locate the responsible network faults.\nIn the 1997 IEEE paper, Fault Isolation and Event Correlation for Integrated Fault Management, the authors, S. Katker and M. Paterok, describe a state-of-the-art algorithm for alarm correlation. The Katker and Paterok algorithm, however, has several disadvantages. First, the algorithm processes alarms very inefficiently. As noted above, a single fault may trigger a large number of network alarms. For example, one fibre cut can result in hundreds of thousands of alarms being reported from circuits supported by the fibre. The Katker and Paterok algorithm initiates a large number of computing threads, each of which ultimately results in the same conclusion. Thus, an excessive amount of time and computational resources are utilized. Additionally, the Katker and Paterok algorithm fails to correlate network element (NE) alarms that are caused by a faulty NE that does not itself generate an alarm.\nIn order to overcome the disadvantages of existing solutions, it would be advantageous to have a system and method of correlating large numbers of network alarms which greatly reduces the time and computational resources utilized, and supports near real-time alarm correlation. The present invention provides such a system and method."}
{"text": "1. Field of the Invention\nThe present invention pertains generally to the field of solid-state X-ray imagers and displays, and more particularly is an improved method that structurally alters the optical path to reduce or avoid radiation damage to the semiconductor components used to process the detected X-ray images.\n2. Background of the Invention\nAs used in this disclosure, X-rays are defined as ionizing electromagnetic radiation that is damaging to semiconductor-based image sensor arrays. X-rays also include the radiation known as “extreme ultraviolet radiation” and “gamma rays”. Since few X-rays with energies exceeding 10 KeV are captured by semiconductor-based image sensor arrays (Si, Ge, etc.), the X-ray energies must be converted into a detectable form. The image sensor arrays are processed on silicon and are only sensitive to light with wavelengths at or near the visible spectrum. Therefore, the arrays require an X-ray-to-visible-light converter in order to detect the X-rays. To this end, X-ray sensitive scintillating materials, such as Gd2O2S:Tb (GOX or GADOX), CsI(Tl) or CdWO4 have been used. These materials greatly enhance the detection efficiency of higher energy X-rays in silicon based sensor arrays through the ability of the scintillating materials to scintillate and emit visible light photons proportional to the X-ray energy. The visible light photons are converted to electrical signals by a silicon based image sensor array, such as a Linear Photodiode Array (PDA). When the image sensor array is read out, the array sequentially produces a stream of electrical video signals from each photo-element with amplitudes proportional to the intensity of the X-ray pattern that impinges on the photo-elements.\nHowever, a problem arises in that the scintillation layer on top of the silicon photo-elements will not absorb the X-ray photons completely. Some portion of the X-ray particles penetrates the scintillation layer and is captured by the image sensor array structure, causing irreversible radiation damage to the image sensor array. Therefore, if the image sensor array used in the X-ray imaging system lies in the X-ray path and is not isolated or protected from X-ray exposure, radiation damage will be inflicted on the silicon image sensor array. As a result, the silicon array used in an X-ray imaging system has a limited useful lifetime.\nFIG. 1 is a simplified electrical block diagram describing the signal processing required for an X-ray detector system. Since the present invention involves only the optical and mechanical structures of such systems, the generalized electrical block diagram shown in FIG. 1 is demonstrative of the signal processing circuitry used in all the systems described herein.\nIn FIG. 1, the detector is a CMOS device with an image sensor array, a PDA, and the readout control circuits for the array. As is known in the art, the PDA is an array of photodiodes with on-chip control circuits for scanning and reading out video signals. In FIG. 1, the PDA is shown with two rows of photodiodes (PD), Row A and Row B. Row A is a dummy row of dark photodiodes used as a reference to differentially cancel any common mode noise from the active video signal, Row B. Row A is covered with metal to shield the photodiodes from light exposure.\nRow B has a light sensing area exposed through a narrow slit in the metal to form a narrow aperture over the length of its read line. When the active photodiodes are exposed to imaging light, each diode collects the photons in the immediate area and converts them to signal charges. The signal charges are stored in the depletion layer capacitance of each individual photodiode. The stored charges are read out during the scanning readout process of the PDA. During one line-scan time, which is known as the integration time, each photodiode goes through an integration process. In the integration process, each photodiode is read out and then reset to its initial condition to start collecting photons and converting them to charges for the following line-scan time. Since the readout is sequential, while the PDA is continuously scanning, each photodiode sequentially goes through the photon collection (integration) process during one line-scan time.\nThe scanning process is initiated by a start pulse, SP. Since the integration time is equal to the line-scan time, the line rate of the video signal is determined by the time required to generate the start pulse, which initiates the scanning of the shift register, SR. As the SR shifts a pulse through its register, two rows of MOS switches, SW, that are in series with the PD are accessed. The pulse from the SR closes two switches. One switch is on the dummy video line, VLD, and the other switch is on the active video line, VLA. As the pulse from the SR accesses the SW, the charges from the accessed PD flow out on to the VLA. The photon converted charges are sent to the signal processing circuit, VP, where the charges are differentially added to the reference charges from the VLD, digitized, and stored for the host computer to perform image processing.\nTo form a Direct Coupled X-ray Detector, a uniform layer of the scintillating material, SCIN, is deposited directly on the sensing areas of the PDA, or a uniform layer of the scintillating material is placed directly on top of the sensing areas of the PDA. The shaded area with diagonal lines in FIG. 1 shows a SCIN layer that has been deposited over the active PD. SCIN is a uniform coated layer that emits photons when its atoms are excited by the impinging beams of the X-ray. The light energies, proportional to the intensity of the X-ray beams, directly expose the active sensing areas of the image sensors and are processed as describe above.\nFIGS. 2a-b show the optical and mechanical components of one of the current art X-ray detector systems commonly used today, a Direct Coupled Detector System. This system is the least complicated in terms of fabrication and applications, and therefore results in lower cost than other systems. The details of the drawing are limited to components relevant to the present invention.\nFIG. 2a shows an isometric view of the components: the image sensors (IS), the test specimen (TS), the exposing X-ray beam (XPXB), etc. FIG. 2b is a sectional view. In FIGS. 2a-b, the PDA; the test specimen under X-ray imaging (TS); the exposing X-ray beam (XPXB); and the scintillation coating (SCIN) on the sensor die are depicted to show the geometrical relationship among the components involved in X-ray testing of the test specimen, TS.\nIn operation, the X-ray source emanates the exposing X-ray beam and exposes the test specimen. The X-ray flux patterns are modulated by the specimen under testing as the flux pattern passes onto the surface of the scintillation coating. Since the scintillation coating is coated directly onto the surface of the image sensor, the converted light energies proportional to the X-ray flux patterns are integrated by the image sensor array as it generates the image video signals.\nThe Direct Coupled Detector System in FIG. 2 shows that the exposing X-ray beam passes through the test specimen, the scintillating layer, and the image sensor array. Accordingly the image sensor array receives that portion of the X-ray flux which is not absorbed by the scintillation layer, causing radiation damage on the silicon sensor. In many applications this radiation exposure is intolerable because it drastically reduces the lifetime of the image sensor array, thereby requiring continual replacement and maintenance of the X-ray imaging system.\nAlthough applying the scintillating layer directly to the image sensor is intolerable for many applications, the primary advantages of the method arise from its simplicity in structure and the close proximity of the scintillating layer to the image sensor array, which improves imaging resolution. Among the advantages of this system is that the detectors are simple to fabricate, i.e., the detectors can be fabricated by simply applying a SCIN coating process to existing image array sensors, such as the PDA. This is a great advantage in applications where a shorter lifetime X-ray detector system is required, for example, in destructive testing where the measuring equipment is also destroyed.\nAnother advantage of a Direct Coupled Detector System arises from the close proximity of the scintillation layer and the photo-element. Since the scintillating coating is in contact with the image plane of the image sensor array, there is little or essentially no space between them. This close proximity gives the detector the ability to retain its optimum resolution and Modulation Transfer Function (MTF).\nAnother advantage of the system, arising from the close proximity of the scintillating layer and the PDA, is the light coupling efficiency, i.e., there is very little light energy loss in the transmission between the scintillating layer and the PDA. Another advantage, which arises from its simple structure, is that the system can be implemented in a small enclosure. The ability to use the system in a small enclosure also allows the system to be designed as a portable unit.\nHowever, there are also several drawbacks to the Direct Coupled Detector System. The system user must tolerate a shorter lifetime for the X-ray detector system in a given application, and the PDA must be continually replaced. Therefore, the Direct Coupled Detector System has a high maintenance cost, and requires a significant amount of down time.\nAnother disadvantage of the Direct Coupled Detector System arises from the noise properties of the PDA. Sensor noise increases with an increasing dose of radiation exposures due to the build-up of undesirable charges in the oxide and silicon interface. Therefore, as the system is used, the noise level increases and the system signal-to-noise ratio decreases. Since noise build-up is a function of radiation exposure, the system performance degrades slowly while the system is in use.\nA third disadvantage of the Direct Coupled Detector System is that the leakage current of the image sensor increases as the interface charge builds up during operation under X-ray exposure. As the leakage current increases, the storage space in the photodiode is decreased until it is rendered unusable. Again since the leakage current build up is proportional to the total dose of X-ray exposure, the system performance will also degrade with time due to leakage current build up.\nA fourth disadvantage in the current art Direct Coupled Detector System is that when some of the X-ray photons pass through the scintillation layer and are absorbed by the photodiode, large signal spikes are created that increase the noise level of the video signal.\nA second prior art system, the Fiber Optics Coupled Detector System, employs a fiber optics bundle to transmit the light from the scintillating layer to the PDA. The object of this system is to isolate the PDA and its electronic components from the exposing X-ray beam. FIGS. 3a-b summarize the optical-mechanical configuration of a Fiber Optics Coupled Detector System. The components of the system are an X-ray source (XS); an exposing X-ray beam (XPXB); a lead shield (LS) with a slit to form an aperture (AP); a test specimen (TS) undergoing X-ray examination; scintillating layer (SCIN) that is coated onto the surface of an optical flat transparent transmission plate (OTP); an image sensor DIP package (PDA); a fiber optic bond (FB); a fiber optic bundle (FO); and an image sensor (IS). The X-ray-to-light converter assembly (ASY) represents the assembly of the scintillating layer, the optical transmission plate, and the fiber optic bond.\nThe X-ray source passes through the aperture to limit the area of the X-ray beam exposure to the neighborhood of the X-ray-to-light converter assembly. The X-ray-to-light converter assembly converts the modulated X-ray flux densities, proportional to the density patterns in the test specimen, to proportional light intensities. The light intensities are coupled into the fiber optic bundle through the fiber bond. The fiber optic bundle couples the light flux down to and through a second fiber bond that couples the light flux onto the surface of the image sensor, where the light flux is integrated and processed. The fiber optic transmission line gives the detector the ability to remotely place the X-ray-to-light converter assembly, hence isolating the PDA and its associated electronic circuits from the path of the exposing X-ray beam. Remotely locating the X-ray detector assembly from the electronic assembly separates the optical path from the X-ray path and achieves the objective of protecting the electronic circuit components from radiation damage.\nA second advantage of the Fiber Optics Coupled Detector System is the preservation of the resolution. A fiber optic bundle has a relatively high optical resolution. However, the fiber optic bundle does create a disadvantage for the system in that glass fiber bundles are expensive and difficult to fabricate.\nA second disadvantage of the Fiber Optics Coupled Detector System is the difficulty of assembly. The glass bundles are difficult to mount and bond. They must be critically aligned and bonded to their transmitting and receiving components to avoid undue optical transmission losses. The alignment constraint is even greater in the case of bonding the fiber optic ends to the surface of the elements of an image sensor because the fiber ends must be cut to exactly match the surface of the image array elements.\nA third disadvantage of the Fiber Optics Coupled Detector System is the constraints imposed on the design of the enclosure. The complicated method used in bonding and mounting the fiber optic bundle requires a supporting structure within the enclosure. The supporting structure, which needs to be flexible enough to make initial adjustments, must also serve as a rigid mount to ensure that the bonded ends remain stationary in transportation and operation. Especially critical in adjusting and mounting are the contacts between the scintillating layer and the optical flat transparent transmission plate and at the fiber bond on the image sensor surfaces. The degree of careful handling required becomes even greater in a two-dimensional application of this X-ray system.\nThere are two inventions that are specifically aimed at eliminating and reducing the disadvantages associated with the above mentioned Direct Coupled and Fiber Optics Coupled X-ray Detector Systems. In U.S. Pat. No. 7,463,717 B2, a rod lens array is used to focus the visible light onto an image sensor array (PDA) after the X-ray flux has been converted as depicted in FIG. 4. The rod lens array and the image sensor array is offset by an acute angle relative to the centerline of the X-ray flux path from the X-ray source, such that electrical components of the detector system are removed from the X-ray flux path.\nIn U.S. Pat. No. 7,463,716 B2, in addition to using a rod lens array for focusing the visible light onto the image sensor array, as shown in FIG. 5, a reflector (a prism or a mirror) is used to reflect the light by ninety (90) degrees and thus remove the electrical components of the detector system from the X-ray flux path. Although both inventions remove the problems associated with direct exposure of X-ray flux onto the image sensor array and other electrical components, the rod lens array has limited light collection efficiency, resulting in loss of light signal. Depending on the models, a rod lens array normally has a light collection angle of about a few tenths of a degree; as a result a significant amount of light signal will be lost.\nAccordingly, it is an object of the present invention to provide an X-ray detector system that is long-lived, compact, low cost, and that has a simple mechanical structure that lends itself to simple production assembly with minimal requirements for alignment, adjustment and calibration testing.\nAnother object of the present invention is to reduce X-ray exposure on components that are sensitive to radiation damage by completely isolating or shielding the components in the detector system from X-ray exposure.\nA further objective of the current invention is to provide an X-ray detector system with high signal light collection efficiency and thus reduce the power requirement of the X-ray source."}
{"text": "Polyphosphoric acid (PPA) is typically made by the polymerization of phosphoric acid via a thermocondensation process. The resulting PPA product has a concentration that can be as high as 118% wt expressed as H3PO4 content. PPA is composed of polymers having different chain lengths, and the composition of PPA varies with its concentration. PPA is used in a multitude of applications where mildly strong acidity and dehydration properties are needed. However, one of the drawbacks of PPA is that, as its concentration increases, its viscosity also increases.\nAs a result of the increased viscosity of PPA, in order to easily use PPA for many processes, PPA must be heated to a sufficient temperature to reduce its viscosity to a level that it can be more easily handled and used, for example by pumping. In some processes, however, increased PPA temperature might not be desirable. In addition, PPA is sometimes stored in plastic containers that are not able to withstand the required higher temperature to allow the PPA to be pumped.\nAccordingly, it would be desirable to have a PPA composition that provides the desired PPA concentration while having a viscosity that allows easy handling at lower temperatures than required by current PPA products."}
{"text": "Self-closing measuring valves, particularly as embodied in lavoratory faucets, have been known in the art for many years. This type faucet is generally characterized by having a manually depressable handle which, when depressed, initiates flow through the faucet. Unless the handle is held in this manually depressed condition, the operating mechanism of the faucet valve acts to slowly close the valve, returning the handle to its uppermost position in readiness for a succeeding manual actuation of the slow self-closing faucet.\nSlow closing lavoratory faucets offer the advantage of assuring that the faucet valve is closed in a matter of a few seconds time so that thereby water usage is conserved and, more importantly, that in the absence of an intending user, the faucet is fully shut off against continued or indefinite water flow through the faucet if the last user intentionally or otherwise overlooked turning off the faucet. Slow closing lavoratory faucets are particularly valuable in conserving water when employed in public washrooms, large institutions and the like.\nThe art also has recognized the desirability in slow self-closing lavoratory faucets of having an adjustability capability for the faucet valve so that the time during which the slow closing operation takes effect after handle depression and manual release of the handle initiates water flow can be changed. In some applications for these slow closing lavoratory faucets, it may indeed be desirable to have a longer flow duration than in installations where a relatively brief flow duration occurs between handle depression and water shut-off.\nFrequently, slow self-closing lavoratory faucets in the prior art have employed a dash-pot form of delay mechanism. This type delay mechanism usually finds it necessary to incorporate a restrictor valve which in effect meters or restricts the flow of fluid into the mechanism, thereby slowing the mechanism down to gain the desired slow closing of the faucet valve which is under control of this mechanism.\nThis sort of restrictor valve is frequently quite sensitive to the presence of dirt or foreign material such as rust particles which are invariably present in large, and particularly old, city water systems. This foreign material flowing through the water system into the restrictor valve of the slow self-closing lavoratory faucet tends to build up and eventually clog the restrictor valve. Thereupon difficult and expensive cleaning of the entire slow self-closing lavoratory faucet frequently becomes necessary.\nAttempts to handle this clogging problem by the use of filters or other screening techniques have only transferred the clogging problem from the restrictor valve becoming clogged to the filter or screen leading to the restrictor valve becoming clogged. In either event the slow self-closing valve requires expensive and frequent cleaning.\nFurther, prior art solutions to obtaining an adjustability capability for varying flow duration have not been entirely satisfactory. The adjustable feature to gain this varying flow duration capability is too often inaccessible from the exterior of the slow closing valve faucet, thereby making it necessary to undertake major disassembly of the faucet simply to obtain variation in the duration of the water flow after the handle is manually depressed."}
{"text": "JP 07-64930 A discloses a microcomputer mutual monitoring method in a CPU backup system in which two CPUs are used as a main CPU and a backup CPU to complement each other. According to this microcomputer mutual monitoring method, if the main CPU becomes a non-operative state due to a failure of the main CPU midway through processing, the backup CPU begins to operate from the process of the task in progress.\nFurther, JP Patent 4003420 discloses a processing apparatus configured to reset a main microcomputer and a sub-microcomputer by stopping a run pulse signal if the sub-microcomputer cannot execute each control process of a calculation monitoring process within a corresponding processing time.\nIn a system including two microcomputers (referred to as a main microcomputer and a sub-microcomputer herein), they monitor mutually each other, and if one microcomputer detects an abnormal event of the other microcomputer, it resets the other microcomputer to perform an attempt to restore it.\nIn general, the main microcomputer is subjected to a run pulse check by a monitoring circuit other than monitoring by the sub-microcomputer, such as a run pulse check, a communication check, or an ALU calculation check; however, since a requirement is too complicated for the monitoring circuit to implement the ALU calculation check, etc., the ALU calculation check, etc., are not performed during the reset of the sub-microcomputer.\nIn particular, in recent years, the level of functional integration of ECUs is increasing for cost reduction, and there is a case where the reset of the sub-microcomputer is desired when control software installed in the sub-microcomputer detects an error event. At that time, even if a condition required to be met to reset the microcomputer is different for the respective installed systems, it is inevitable to reset the sub-microcomputer as a whole if any one installed system needs reset, because it is not possible to reset only a part of the sub-microcomputer due to a microcomputer configuration.\nAccording to a configuration in which the processing of the main microcomputer continues in spite of not being capable of monitoring the main microcomputer during the reset of the sub-microcomputer, there is a problem that reliability as a system is reduced.\nOn the other hand, such a configuration may be contemplated in which the processing of the main microcomputer discontinues during the reset of the sub-microcomputer, considering that it is not possible to monitor the main microcomputer. However, according to such a configuration, since the main microcomputer is reset regardless of whether the main microcomputer is abnormal or normal, there is a problem that marketability is reduced."}
{"text": "This application relates to a computer network and, more specifically, to a method and apparatus for allowing both high-speed and regular-speed access to a computer network.\nThe Internet is an example of a TCP/IP network. The Internet has over 10 million users. Conventionally, access to the Internet is achieved using a slow, inexpensive method, such as a terrestrial dial-up modem using a protocol such as SLIP (Serial Line IP), PPP, or by using a fast, more expensive method, such as a switched 56 Kbps, frame relay, ISDN (Integrated Services Digital Network), or T1.\nUsers generally want to receive (download) large amounts of data from networks such as the Internet. Thus, it is desirable to have a one-way link that is used only for downloading information from the network. A typical user will receive much more data from the network than he sends. Thus, it is desirable that the one-way link be able to carry large amounts of data very quickly. What is needed is a high bandwidth one-way link that is used only for downloading information, while using a slower one-way link to send data into the network.\nCurrently, not all users have access to high speed links to networks. Because it will take a long time to connect all users to networks such as the Internet via physical high-speed lines, such as fiber optics lines, it is desirable to implement some type of high-speed line that uses the existing infrastructure.\nCertain types of fast network links have long propagation delays. For example, a link may be transmitting information at 10 Mbps, but it may take hundreds of milliseconds for a given piece of information to travel between a source and a destination on the network. In addition, for even fast low-density links, a slow speed return-link may increase the round trip propagation time, and thus limit throughput. The TCP/IP protocol, as commonly implemented, is not designed to operate over fast links with long propagation delays. Thus, it is desirable to take the propagation delay into account when sending information over such a link.\nThe present invention overcomes the problems and disadvantages of the prior art by allowing a user to download data using a fast one-way satellite link, while using a conventional low-speed Internet connection for data being sent into the network. The invention uses a xe2x80x9cspoofingxe2x80x9d technique to solve the problem of the long propagation delays inherent in satellite communication.\nIn accordance with the purpose of the invention, as embodied and broadly described herein, the invention is a network system that forms a part of a network, comprising: a source computer, having a link to the network; a destination computer, having a link to the network; a satellite interface between the source computer and the destination computer, wherein information passes from the source computer to the destination computer; means in the destination computer for requesting information from the source computer over the network; means for receiving an information packet sent from the source computer in response to the request and for sending the information packet to the destination computer over the satellite interface; and means for sending an ACK message to the source computer in response to receipt of the information packet, wherein the ACK message appears to the source computer to have come from the destination computer.\nIn further accordance with the purpose of the invention, as embodied and broadly described herein, the invention is a gateway in a network system that forms a part of a TCP/IP network, wherein the network includes a source computer having a link to the TCP/IP network and a link to a high speed satellite interface, and a destination computer having a link to the TCP/IP network and a link to the high speed satellite interface, the gateway comprising: means for receiving an information packet sent from the source computer and for sending the information packet to the destination computer over the satellite interface; and means for sending an ACK message to the source computer in response to receipt of the information packet, wherein the ACK message appears to the source computer to have come from the destination computer.\nObjects and advantages of the invention will be set forth in part in the description which follows and in part will be obvious from the description or may be learned by practice of the invention. The objects and advantages of the invention will be realized and attained by means of the elements and combinations particularly pointed out in the appended claims."}
{"text": "A dielectric antenna such as a dielectric rod antenna is a surface-wave antenna in which an end-fire radiation pattern is produced by propagation of a surface wave on a tapered dielectric rod. Dielectric rod antennas provide significant performance advantages and are low cost alternatives to free space high-gain antennas at millimeter-wave frequencies and the higher end frequencies of the microwave band. Conventional dielectric antennas required for radar level gauging do not withstand high temperatures. Additionally, such antennas must be installed via small tank nozzles which can affect gain, return loss, and side lobes over the radar bandwidth issues.\nBased on the foregoing, it is believed that a need exists for an improved dielectric hollow antenna, which will be described in greater detail herein."}
{"text": "This invention relates to a fluid drive oil recovery process which utilizes an injection of CO.sub.2, surfactant and water into subterranean oil reservoir. More particularly, the invention relates to such a process in which the surfactant selected for use is a particular member of a relatively highly chemically stable and salt-tolerant class of surfactants and is uniquely suited for use in the reservoir to be treated.\nNumerous patents have been issued on materials and techniques which are pertinent to an oil recovery process that utilizes an injection of CO.sub.2, surfactant and water. The U.S. Pat. Nos. 2,226,119; 2,233,381 and 2,233,382 describe polyalkoxylated alcoholic or phenolic surfactants which are generally useful in aqueous liquid fluid drive oil recovery processes. U.S. Pat. No. 2,623,596 indicates that an increased oil recovery may be obtained by a fluid drive process which injects highly pressurized CO.sub.2. U.S. Pat. No. 3,065,790 indicates that, in a fluid drive process, the cost effectiveness of highly pressurized CO.sub.2 may be increased by injecting a slug of the CO.sub.2 ahead of a cheaper drive fluid. U.S. Pat. No. 3,330,346 indicates that almost any process for forming foam within a reservoir may be improved by using as the surfactant a polyalkoxylated alcohol sulfate of an alcohol containing 10 to 16 carbon atoms. U.S. Pat. No. 3,342,256 indicates that, in a fluid drive process, the oil-displacing efficiency of a CO.sub.2 slug may be increased by including water and a foaming surfactant within that slug. U.S. Pat. No. 3,529,668 indicates that, in a fluid drive process, the efficiency of a slug of foamed CO.sub.2 may be increased by displacing it with specifically proportioned alternating slugs of gas and liquid. U.S. Pat. No. 4,088,190 indicates that, in a fluid drive process, the heat stability and durability of a CO.sub.2 foam may be increased by using an alkyl sulfoacetate surfactant. U.S. Pat. No. 4,113,011 indicates that in a CO.sub.2 foam drive, the problems of low salt tolerance with are typical of both the surface-active sulfates of polyalkoxylated alcohols containing 10 to 16 carbon atoms recommended by U.S. Pat. No. 3,330,346 and the alkyl sulfoacetate surfactants recommended by U.S. Pat. No. 4,088,190 may be avoided by using a surfactant sulfate of a polyalkoxy alcohol containing only 8 or 9 carbon atoms and injecting that surfactant ahead of the CO.sub.2."}
{"text": "A vending machine generally refers to a machine that dispenses items after the customer inserts currency or credit into the machine.\nElements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions and/or relative positioning of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present invention. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments of the present invention. Certain actions and/or steps may be described or depicted in a particular order of occurrence while those skilled in the art will understand that such specificity with respect to sequence is not actually required. The terms and expressions used herein have the ordinary technical meaning as is accorded to such terms and expressions by persons skilled in the technical field as set forth above except where different specific meanings have otherwise been set forth herein."}
{"text": "Non-overlap circuits are typically integrated in ICs that frequently switch. For example, non-overlap circuits may be implemented in half-bridge drivers and alternating current (“AC”) to direct current (“DC”) converters. These non-overlap circuits provide signals to switching transistors, which may be implemented as metal oxide semiconductor field effect transistors (“MOSFETs”), so the transistors do not overlap. Put another way, the signals provided by the non-overlap circuits prevent the low-side and high-side power switches from switching on at the same time.\nIn high-voltage operations, parasitic capacitance of the switching device is large and varies with process. Timing controls circuits are used to adapt the non-overlap circuits to avoid shoot-through currents that may damage the switching devices."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic absolute encoder that uses rotation detectors each including a gear fixed to a rotary shaft and having a predetermined number of teeth and at least one magnetic detecting element for detecting magnetic flux passing through the teeth to output an electrical signal relating to the position of the teeth.\n2. Related Art\nJapanese Patent Application Publication No. 11-237256 (JP11-237256A) discloses a rotation detecting device that can be utilized as a magnetic incremental encoder. Japanese Patent Application Publication No. 2008-180698 (JP2008-180698A) discloses a magnetic absolute encoder that utilizes the basic structure of a magnetic incremental encoder."}
{"text": "A variety of tile products are widely used in the building industry. Such products range from inexpensive types of products, such as vinyl-asbestos flooring tile, to substantially more costly tile products, such as marble tile. Although the latter are preferred due to their appearance, the use of such materials is not without problems. Ceramic and marble tiles are brittle and tend to crack; therefore, they require special installation techniques involving the setting of the tiles into ultimately rigid cement/mortar bases on a solid dimensionally stable subsurface, and the subsequent filling of the joints between the tiles with grout. Such procedures greatly add to the expense and are not always satisfactory."}
{"text": "1. Field of the Invention\nThe present invention relates to wireless communication, and, more particularly, to a method and apparatus for performing a hybrid automatic repeat request (HARQ) in a wireless communication system based on Time Division Duplex (TDD).\n2. Related Art\nLong Term Evolution (LTE) based on 3rd Generation Partnership Project (3GPP) Technical Specification (TS) Release 8 is the leading next-generation mobile communication standard.\nAs disclosed in 3GPP TS 36.211 V8.7.0 (2009-05) “Evolved Universal Terrestrial Radio Access (E-UTRA); Physical Channels and Modulation (Release 8)”, in LTE, a physical channel can be divided into a Physical Downlink Shared Channel (PDSCH) and a Physical Downlink Control Channel (PDCCH), that is, downlink channels, and a Physical Uplink Shared Channel (PUSCH) and a Physical Uplink Control Channel (PUCCH), that is, uplink channels.\nA PUCCH is an uplink control channel used to send uplink control information, such as a Hybrid Automatic Repeat reQuest (HARQ), an acknowledgement/not-acknowledgement (ACK/NACK) signal, a Channel Quality Indicator (CQI), and a Scheduling Request (SR).\nMeanwhile, 3GPP LTE-Advanced (A) that is the evolution of 3GPP LTE is in progress. Technology introduced into 3GPP LTE-A includes a carrier aggregation.\nA carrier aggregation uses a plurality of component carriers. A component carrier is defined by the center frequency and a bandwidth. One downlink component carrier or a pair of an uplink component carrier and a downlink component carrier correspond to one cell. It can be said that a terminal being served using a plurality of downlink component carriers is being served from a plurality of serving cells.\nIn a Time Division Duplex (TDD) system, the same frequency is used in uplink and downlink. Accordingly, one or more DL subframes are associated with an UL subframe. The “association” means that transmission/reception in the DL subframe is associated with transmission/reception in the UL subframe. For example, when a transport block is received in a plurality of DL subframes, a terminal sends HARQ ACK/NACK (hereinafter referred to as ACK/NACK) for the transport block in an UL subframe associated with a plurality of DL subframes. Here, a minimum time is necessary to send the ACK/NACK. This is because the time taken to process the transport block and the time taken to process the ACK/NACK are necessary.\nMeanwhile, a plurality of serving cells can be introduced into a TDD system. That is, a plurality of serving cells can be assigned to a terminal. In this case, in the prior art, it was assumed that the same uplink-downlink (UL-DL) configuration was used in all the serving cells. The UL-DL configuration is information indicating whether each subframe within a radio frame used in TDD is an UL subframe or a DL subframe. In the next-generation wireless communication system, however, to use different UL-DL configurations in serving cells is also taken into consideration. In this case, how an HARQ will be performed using what method is problematic."}
{"text": "It is known that a luminescent aromatic molecule embedded in plastic is subject to quenching by oxygen present in the gas or liquid in contact with the plastic. This phenomenon was reported by Bergman (Nature 218:396, 1966), and a study of oxygen diffusion in plastic was reported by Shaw (Trans. Faraday Soc. 63:2181-2189, 1967). Stevens, in U.S. Pat. No. 3,612,866, ratios the luminescence intensities from luminescent materials dispersed in oxygen-permeable and oxygen-impermeable plastic films to determine oxygen concentration. Lubbers et al. in U.S. Pat. No. 4,003,707 proposed the possibility of positioning the emitting substance at the end of an optical fiber. Peterson et al. in U.S. Pat. No. 4,476,870 also employs the quenching of an emitting molecule in plastic at the end of an optical fiber. Both Lubbers and Peterson reference emission against scattered exciting light.\nThe quenching of the luminescence of an emitter at the end of an optical fiber has been used in temperature sensors. For temperature probes the emitters are generally solid phosphors rather than an aromatic molecule embedded in plastic, since access by molecules from the environment is not desirable. Various methods have been used to measure the amount of quenching: (i) Quick et al. in U.S. Pat. No. 4,223,226 ratios the intensity at one wavelength of the emission against another; (ii) Quick et al. also proposes determining the length of time it takes for the signal to fall from one level to another; (iii) Samulski in U.S. Pat. No. 4,245,507 (reissued as U.S. Pat. No. Re. 31,832) proposes to measure quenching by determining the phase of the emitted life. In a very recent patent for temperature sensing at the end of an optical fiber, Hirschfeld in U.S. Pat. No. 4,542,987 proposes, in addition to method (i), that (iv) emission lifetime be used to measure quenching and that (v) Raman scattered light can be used as a reference.\nEastwood and Gouterman (1970) noted generally with respect to Pd and Pt porphyrin complexes that their \"relatively high [emission] yields and short triplet lifetimes . . . may make these systems useful as . . . biological probes for the presence of oxygen.\" More recently, Bacon and Demas in UK Patent Application No. 2,132,348A propose the use of, inter alia, porphyrin complexes of VO.sup.2+, Cu.sup.2+, Pt.sup.2+, Zn.sup.2+ and Pd.sup.2+ or dimeric Rh, Pt, or Ir complexes for monitoring oxygen concentration by emission quenching of intensity or lifetime. Suitable ligands would reportedly be etioporphyrin, octaethylporphin, and porphin."}
{"text": "1. Field of the Invention\nThe present invention relates to a dust control composition for fertilizer for the purpose of reducing the dust levels present in the fertilizer, reducing subsequent dust formation, and to reduce the tendency of the fertilizer particles to agglomerate or cake during storage and transportation.\n2. Description of Related Art\nMethods for the manufacture of fertilizers (inorganic, organic, or micronutrient) as well as methods for processing these fertilizers into particles via prilling, granulating, compaction or other techniques are well known. The resulting fertilizers often contain an undesirable level of particles fine enough to become airborne dust. This dust is produced during the manufacture, storage and transportation of the fertilizer particles. The dust can be the result of mechanical abrasion encountered during movement of the fertilizer particles, continued chemical reactions or curing processes after the initial particle formation, the action of moisture migration through the fertilizer during storage or ambient temperature and humidity conditions.\nFertilizer dust dissemination poses safety, health, environmental, housekeeping and maintenance problems for fertilizer producers, distributors and consumers. For instance fertilizer dust has raised health concerns due to human and animal inhalation thereof. It is also a concern when fertilizer dust becomes airborne which can lead to the loss of agronomic and economic value, while potentially contributing to the contamination of surface water ecosystems.\nThe use of oils, waxes, blends of oil and wax, and emulsions based on these products have been known for a long time. These oils and waxes can be petroleum or vegetable based. For instance, in 1977 Frick suggested that petroleum based products be used to control dust from agricultural fertilizers, See xe2x80x9cPetroleum Based DCA\"\"s to Control Fugitive Dustxe2x80x9d, Frick, Proceedings of the Annual Meeting of the Fertilizer Industry Round Table, Series 27, pages 94-96. However there are disadvantages involved in using these treatment methods. Over time oils tend to volatilize and/or be adsorbed into the fertilizer particle and loose their effectiveness. Waxes are also ineffective and difficult to handle because they absorb into the fertilizer particle at temperatures above their melt point and do not spread or coat the fertilizer particle surface at temperatures below their melt point. In addition, both oils and waxes have limited binding properties that are essential for long term fertilizer dust control.\nOther proposed dust control methods include application of other liquids such as lignosulfonate solutions, molasses solutions, urea solutions, mixtures of these solutions, other fertilizer solutions, amines, surfactants, polymers and even water. Examples are U.S. Pat. No. 5,360,465 to Buckholtz et al. and U.S. Pat. No. 5,328,497 to Hazlett. These methods have a number of disadvantages as well. Due to the water present, aqueous solutions and emulsions tend to accelerate the formation of fertilizer dust and exacerbate the fertilizer particles caking tendencies. These treatments also tend to loose their binding properties as the solutions and emulsions dry, thereby becoming ineffective as long term dust control agents.\nEuropean Patent 0320987 discloses the use of a conditioning agent comprising 10-60% wax, 30-90% oil and 0.3-10% by weight of a high-molecular weight viscoelastic elastomer such as polyisobutylene. U.S. Pat. No. 5,603,745 to Pettersen et al. discloses the use of a conditioning agent comprising 10-50% wax, 40-90% oil and 1-30% of a oil soluble and wax miscible resin such as esters of polymerized resin, esters of stabilized resin acids or non-crystallized tall oil resin. While these conditioning agents provided an improvement in dust control over oils, waxes, and oil/wax blends, they do not provide the degree of binding required for effective long-term dust control. This is because the majority of these conditioning agents still consist of poor performing oils and waxes.\nThe vast majority of commercially produced fertilizers are treated with a conditioning agent of some type to reduce dust levels. For instance, in 1995 Ogzewalla suggested several characteristics needed for an effective dust control conditioning agent in xe2x80x9cFertilizer Dust and Dust Control Coating Agentsxe2x80x9d, Ogzewalla, Proceedings of the Annual Meeting of the Fertilizer Round Table, 45th Meeting, 1995, pages 95-100. These characteristics included the ability to bind dust back to the surface of the fertilizer granule, resist absorption into the fertilizer granule surface, and the ability to spread or coat the fertilizer granule surface. The example given in this article shows that several commercially available conditioning agents were able to reduce dust levels in a diammonium phosphate fertilizer from 600 ppm to between 170-70 ppm. In addition, the xe2x80x9cSixth Annual 1996 Granular Fertilizer Surveyxe2x80x9d, 1996, ARR-MAZ Products, L.P., shows the dust levels found in twenty two samples of diammonium phosphate collected from fertilizer production facilities across North America. The producers, during manufacture, had treated all of these fertilizer samples with a conditioning agent. Dust levels were shown to range from 425 ppm to 55 ppm and averaged 125 ppm.\nOne of the main objectives of the invention was to develop a superior agricultural composition comprising of fertilizer particles having low dust levels and a reduced tendency to cake during the long term storage and handling conditions normally encountered by commercial fertilizer products. Another objective was to develop a superior conditioning agent that is fluid at application temperatures and can be applied by conventional coating or conditioning equipment. A further objective was to arrive at a conditioning agent that would not effect the fertilizers handling characteristics or flowability.\nThe main problem to be solved was to obtain a superior conditioning or dust control agent that is fluid and flexible enough to spread over the surface of the fertilizer granules during the coating process, and yet still had enough binding properties to adhere ambient dust to the surface of the granule and reduce dust formation during subsequent storage and handling. Further, it was important that the resulting treated fertilizer granules can not become too sticky for handling by conventional means, even when treated at relatively high application rates. In addition, the resultant conditioning agent should be easy to apply on the fertilizer particles and be non-toxic to the soil and plants. This latter requirement implies that the various components must be environmentally acceptable. It is also desirable that the conditioning of the fertilizer particles be accomplished in one step with the required protection obtained during this step. A further requirement was that the treated fertilizer be completely soluble a few days after to the soil, and that the conditioning agent be degradable in the soil.\nIn view of the above stated requirements, an investigation was started for an improved conditioning or dust control agent by studying ways to improve the properties of the oils and waxes used as major components of the most common conditioning agents. Both petroleum and vegetable based oils and waxes have some value in these conditioning agent formulations. The effectiveness of any specific oil was found to be determined by the oils physical properties, in particular the combination of viscosity and tackiness. Accordingly, this component, the oil/wax, could be substituted with a new component, another oil/wax, having somewhat different properties. Further investigation showed that selection of optimal components could be of importance.\nThe search for optimal components resulted in investigation of methods that could change or modify these components. The selection of a vegetable oil or wax and then subjecting this vegetable oil or wax to an oxidation process was found to have a measurable effect on reducing fertilizer dust and dust formation.\nTypical examples of materials which can be oxidized and used as starting points for the oxidized component include, but are not limited to, crude oils such as corn oil, canola oil, cottonseed oil, sunflower oil, soy oil, linseed oil, castor oil, and tall oil as well as their distillation products and distillation residues such as distilled tall oil, tall oil pitch and tall oil bottoms or mixtures thereof.\nThe compounding or mixing of lower viscosity oils or waxes into the oxidized oil are useful, but not required, to standardize the physical properties of the finished conditioning agent such as viscosity, melt point and tackiness. This insures that the conditioning agent has the desired properties at application temperatures which provide for improved dust control abilities and also allows for application by conventional coating or conditioning equipment.\nNon-oxidized oils useful for compounding with the oxidized oils include, but are not limited to, white oil, refined mineral oils, and vegetable oils such as corn oil, canola oil, cottonseed oil, sunflower oil, soy oil, linseed oil, castor oil and tall oil. Oils having moderate viscosity, low volatility, and high flash point are preferred.\nWaxes useful for compounding with the oxidized oils include but are not limited to intermediate waxes, paraffin waxes, micro-crystalline waxes, carnauba wax, vegetable waxes and mixtures thereof. Waxes with low congealing points are preferred for low temperature applications, while waxes with high congealing points can be useful in high temperature applications.\nThe oxidized oil component of the novel conditioning agent must be soluble or miscible with the oil and/or wax components. Further, it must result in a coating with the viscosity and tackiness required for dust control while maintaining the fluidity required ease of application. The resulting coating should control ambient dust levels, reduce dust formation and reduce caking tendencies without adversely effecting the fertilizers handling characteristics. These components should also be environmentally acceptable and degrade in the soil as explained above. Within this framework, it was then found that effective relative amounts in weight % of these components should be:\nOxidized Oil: 10-100%, preferably 30-90% typically\nOil: 0-90%, preferably 10-70% typically\nWax: 0-90%, preferably 10-70% typically\nThe present invention is therefore a dust control composition for fertilizer particles, the composition comprising 10-100% by weight oxidized oil, 0-90% by weight non-oxidized oil, and 0-90% by weight wax. The fertilizer particle is selected from the group consisting of ammonium phosphate, potash, granulated single super phosphate, granular triple super phosphate, NP-fertilizer and NPK-fertilizer.\nIn a typical preferred application, the composition comprises 30-90% by weight oxidized oil, 10-70% by weight non-oxidized oil, and 10-70% by weight wax.\nThe oxidized oil is selected from the group consisting of corn oil, canola oil, cottonseed oil, sunflower oil, soy oil, linseed oil, castor oil, tall oil, mixtures thereof, and distillation products and distillation residues thereof; the non-oxidized oil is selected from the group consisting of white oil, refined mineral oil, vegetable oil, and mixtures thereof; and the wax is selected from the group consisting of intermediate waxes, paraffin waxes, micro-crystalline waxes, carnauba waxes, vegetable waxes and mixtures thereof.\nThe present invention is also a method of controlling dust from fertilizer particles which comprises treating the fertilizer particles with a composition as described above."}
{"text": "This invention relates generally to film cooled combustor liners for use in gas turbine engines and more particularly to such combustor liners having regions with closely spaced cooling holes.\nA gas turbine engine includes a compressor that provides pressurized air to a combustor wherein the air is mixed with fuel and ignited for generating hot combustion gases. The fuel is injected into the combustor through fuel tubes located at uniformly spaced injection points around the combustor. These gases flow downstream to one or more turbines that extract energy therefrom to power the compressor and provide useful work such as powering an aircraft in flight. Combustors used in aircraft engines typically include inner and outer combustor liners to protect the combustor and surrounding engine components from the intense heat generated by the combustion process. A variety of approaches have been proposed to cool combustor liners so as to allow the liners to withstand greater combustion temperatures. One such approach is multi-hole film cooling wherein a thin layer of cooling air is provided along the combustion side of the liners by an array of very small cooling holes formed through the liners. Multi-hole film cooling reduces the overall thermal load on the liners because the mass flow through the cooling holes dilutes the hot combustion gas next to the liner surfaces, and the flow through the holes provides convective cooling of the liner walls.\nIn the assembled combustor, certain portions of the combustor liners are aligned with the injection points defined by the circumferential location of the center of the fuel tubes. These locations are hereinafter referred to as xe2x80x9ccup centersxe2x80x9d. In operation, the flow of combustion gases past these circumferential locations create xe2x80x9chot streaksxe2x80x9d of locally increased material temperatures. The portions of the combustor liners subject to hot streaks can exhibit oxidation, corrosion and low cycle fatigue (LCF) failures after return from field use.\nAccordingly, there is a need for a combustor liner in which cooling film effectiveness is increased in the areas of the liner that are subject to unusually high temperatures and resulting material distress.\nThe above-mentioned need is met by the present invention, which provides a gas turbine combustor liner made up of a shell having cooling holes formed therein, a group of which are disposed upstream of the dilution holes and divided into two sub-groups. The second sub-group of this group of cooling holes is located in circumferential alignment with a hot streak and are more closely spaced than the cooling holes of the first sub-group. The shell may also have additional cooling hole groups disposed between dilution holes in the liner. The additional groups are arranged so as to provide a converging flow in the circumferential direction to provide enhanced cooling to the area of the liner downstream of the dilution holes."}
{"text": "The release carriers of the prior art are made by applying a polyolefin, preferably polypropylene, onto a continuous sheet of release paper, passing the polypropylene covered paper around a chill roll and rolling the release carrier onto a core for further processing, storage, shipment and handling. Since the polypropylene is set while in contact with the chill roll, it takes on the inverted image of the roll. Therefore, it is critical that the surface of the chill roll be defect free.\nEven if the chill roll is maintained defect free, dull streaks are created in the release carrier due to gauge bands (the difference in caliber or thickness) of the release paper in the machine direction. It is common for 3,000 laps of release carrier to be wound into a 42 inch diameter roll. Due to the cumulative effect of the gauge bands, the difference in diameter across the roll is typically as much as 1/4 inch. This leads to rubbing of the backside of the release paper against the adjacent polypropylene surface causing dull streaks in the high gloss polypropylene surface during handling and storage.\nPresently, the surface defects are removed during manufacture of surface coverings such as vinyl floor coverings by pretreating the release carrier at the floor covering manufacturing site. As the release carrier is unrolled, the polypropylene is heated to a temperature greater than its melting point (in excess of 350.degree. F.) and planished. While the dull streaks are removed by this process, the planisher must be maintained defect free.\nFurther, no known prior art release carrier has been made having two different predetermined gloss levels in predetermined areas of the release carrier.\nDavidson, U.S. Pat. No. 3,507,733, discloses the use of a polypropylene coated release carrier in the manufacture of an embossed decorative surface covering.\nErb et al., U.S. Pat. No. 3,773,545, discloses a process for controlling the surface gloss of a vinyl coated floor covering. In particular, they disclosed the use of a hot polishing roll on a cool vinyl coated substrate to improve the gloss.\nO'Sullivan, U.S. Pat. No. 4,478,663, discloses the use of a highly polished chill roll which has depressions of an average depth of about 5 microns and an average area for each depression of less than 16,000 square microns.\nAs evidenced by O'Sullivan, the teachings of the present invention can also be applied to products composed of laminates of plastics and other materials used in a variety of packages and containers such as plastic bags."}
{"text": "In a telecommunications system where a central telecommunications station supports a plurality of subscribers and a controller is provided for controlling one or more such central telecommunications stations, it is necessary to pass control and other messages between the controller and the central telecommunications station. The messages should be handled in a reliable and efficient manner. It should also be possible to detect messages which are lost and to re-send those messages."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for surface treatment.\n2. Description of the Prior Art\nIn recent years, semiconductor devices, quartz oscillators, and the like are manufactured through processing carried out using photolithography technology or the like, that is, so-called micromachining technology is actively employed. By employing micromachining technology, it is possible to directly process substrates. For example, after one surface of a substrate is subjected to processing, the thickness of the substrate can be reduced or holes can be provided in the substrate by subjecting the other surface of the substrate, which is opposite to the processed surface, to etching or the like.\nIn such a case where one surface of a substrate is processed (that is, in a case where one surface of a substrate is subjected to etching), there is a necessity to protect the other surface that has been processed (hereinafter, referred to as a “processed surface”) from an etchant. From such a viewpoint, an apparatus for etching capable of subjecting only a surface opposite to a processed surface (hereinafter, referred to as a “surface to be etched”) to etching while protecting the processed surface has been developed (see Japanese Patent Laid-open No. Hei 7-111257, for example). The apparatus for etching disclosed in Japanese Patent Laid-open No. Hei 7-111257 includes a supporting jig and a pressing jig to be detachably attached to the supporting jig. When the pressing jig is attached to the supporting jig with a substrate being placed on the supporting jig, the pressing jig presses the peripheral portion of the substrate. In this way, the substrate is secured to the apparatus for etching.\nFurther, the apparatus for etching has O-rings provided at predetermined positions. While the substrate is being secured to the apparatus for etching, the O-rings are in close contact with a surface to be etched of the substrate. This makes it possible to prevent an etchant from being reached to a processed surface of the substrate. Therefore, even when the apparatus for etching to which the substrate is secured is immersed in an etchant, only the surface to be etched is subjected to etching.\nHowever, in the case where such an apparatus for etching is used, the outer peripheral portion of the surface to be etched which is pressed by the pressing jig is not exposed to the etchant so that the outer peripheral portion is not treated.\nAs a result, a difference in level is developed on the etched surface at the boundary between a portion that has been in contact with the pressing jig (that is, an untreated portion) and a portion that has not been in contact with the pressing jig (that is, a treated portion). Such a difference in level causes a problem that it is difficult to evenly carry out various treatments for each area of the substrate in post-steps coming after etching."}
{"text": "Field of the Invention\nThe present invention relates to the field of context aware computing and more particularly to contextual mediation.\nDescription of the Related Art\nContext aware computing refers to the treatment of the user environment in system behavior. Aspects of context can be wide-ranging and can include computing device and undertaken task. Context can be used to modify data in a computing system in order most appropriately process the data. Examples include adjusting the presentation of data according to context, or adjusting the responsiveness of logic to events according to the context of processed data. Contextual mediation generally relates to the consideration of the contemporaneous context of a transaction in processing the transaction itself.\nGenerally speaking, the notion of mediation has been proposed as the principle means of resolving semantic inter-operation issues in a computing environment. Mediation architectures often are based upon mediator or wrappers paradigms wherein information flowing from a source can be wrapped into logical views such that the interface in front of each view appears uniform. The logical views, however, can be assembled via a domain-ontology such as an integrated global schema. The mediator or wrapper, in turn, can serve as an intermediary or interface between observing end users and the source of information.\nOnce such source of information can include event data in a computing environment. Event data often can be generated natively or via adapters, through common base event generation. Alternatively, event data can be generated via crawling or extensive indexing of content in the computing system. In a given process in the computing system, however, it can be important to correctly identify relevant data according to a provided context such that the identified relevant data accordingly can be consumed by an appropriate consumer."}
{"text": "Surgical operations, including more complex operations where a substantial amount of bleeding may occur, may require transfusions during the course of the surgery to maintain a sufficient blood volume and blood pressure. Since many blood-borne diseases may exist including hepatitis, cancer and HIV, it is desirable to not require transfusion from another person. Also, if blood or blood components from the same person can be used, the necessity to match blood factors can be eliminated.\nThese disadvantages of receiving transfusions from donors are overcome by self-donation prior to operations. However, operations involving transfusions are not always identified in advance and few patients take the time and effort to go through the procedure. Additionally, a patient may be weakened by removal of blood prior to an operation.\nAutotransfusion, whereby blood retrieved from the patient during the operation is separated so that reusable portions can be reinserted into the patient, is an effective method of overcoming the problems with transfusions. Various autotransfusion type systems currently exist but are somewhat complex to operate. For example, some autotransfusion systems require the operator to memorize a series of system steps to insure that the operator performs operations in the proper sequential order. Failure to perform the step or to perform the step in the proper sequence may cause the system to shut down or may cause morbidity in the patient.\nAdditionally, it is highly useful to have a blood separation system that can efficiently separate platelet and plasma from waste products in the blood. A high degree of efficiency in obtaining platelets has not been previously achieved.\nIt is therefore desirable to have a blood separation system that is highly efficient in extracting platelets from the blood, extracting waste products from the blood, allowing performance of operations in a simple and easy manner that does not require extensive knowledge of the system and processes, and preventing inadvertent or accidental operation of the blood separation device."}
{"text": "The present invention relates generally to polyamides, and more particularly to novel polyamides which exhibit intumescent flame retardant properties, to a process for preparing the polyamides, and to flame retardant coating formulations of the intumescent type containing the novel polyamides.\nThe importance of protecting building materials and other heat and fire-vulnerable substrates against the effects of high temperatures and flames is widely recognized. The use of flame retardant coating formulations to impart flame retardance to these substrates has been known for several years. A class of these flame retardant coating formulations is formulations of the intumescent type.\nIntumescence is a state of being tumid or inflated. An intumescent coating is one that will enlarge or expand to form a cellular structure when exposed to sufficient heat. Coatings of the intumescent type provide protection to heat and/or fire-vulnerable substrates by forming a flame retardant, thermally insulating barrier over the substrate.\nThe flame retardant coating formulations generally employed in the art are multicomponent materials containing an inorganic flame retardant component, (such as a phosphorous containing compound), a suitable film-forming binder, a dispersing agent, fillers and pigments, and an intumescent component. The intumescent component generally contains two ingredients: (1) a nonresinous chemical material, called a carbonific, which forms a large volume of carbonaceous char, and (2) a chemical material, sometimes referred to as a spumific, which upon thermal decomposition releases large quantities of nonflammable gaseous products. Typical flame retardant formulations of the above type are taught in U.S. Pat. Nos. 3,396,129; 3,440,201; 3,449,161; and 3,562,197.\nU.S. Pat. No. 3,297,754 discloses a process for preparing polyamides from .alpha.,.alpha.-disubstituted .beta.-halopropionic acid amides by way of .beta.-lactams. The polymerization proceeds by a ring opening mechanism and consequently results in polyamide compositions differing in structure from the presently claimed compositions. Moreover, there is no teaching that the polyamides produced by the process of U.S. Pat. No. 3,297,754 exhibit intumescent flame retardant properties."}
{"text": "Illumination apparatuses including semiconductor light-emitting devices such as laser emitting diodes (LED) as light sources, instead of fluorescent or incandescent lamps, have been used indoors, or used in buildings or in houses. For example, an illumination apparatus including a light-emitting device is used as a light source for indoor visual inspection of painted surfaces of products, such as home electric appliances and automobiles.\nA semiconductor light-emitting device emits light with a narrow region of wavelengths, and may only emit monochromatic light. To produce white light as illumination light, a plurality of semiconductor light-emitting devices that emit light with different wavelength regions are prepared, and a plurality of light beams with different colors emitted from such semiconductor light-emitting devices are combined to produce white light. A plurality of phosphors that emit light with different wavelength regions using excitation light with the same wavelength are prepared, and a plurality of fluorescence beams with different colors are combined into white light.\nThis method of combining colors allows a light source to produce white light and also light with other spectra for intended use.\nThe light-emitting apparatus described in Japanese Unexamined Patent Application Publication No. 2015-126160 includes two blue light-emitting devices with different peak wavelengths, a green phosphor that is excited by light emitted from the blue light-emitting device to emit green light, and a red phosphor that emits red light to improve color rendering."}
{"text": "This invention pertains to improved polyester moldings and more particularly to those having improved surface characteristics.\nA technical improvement that has made a significant contribution to commercial polyester molding technology is the use of low profile additives to reduce shrinkage during the curing reaction, and to thereby improve dimensional stability and surface smoothness. Low profile additives are thermoplastic polymers such as vinyl acetate polymers, polystyrene, acrylic polymers, and polycaprolactones. There are a number of theories that seek to explain the low profile or anti-shrinkage action of these polymers, but the one that seems to best explain the phenomenon is the following:\nThe low profile additive is at least partly soluble in the uncured polyester/styrene solution. As the polyester/styrene mixture crosslinks, the thermoplastic polymer becomes incompatible or less soluble and at least partly comes out of solution. This action causes a volume expansion that compensates for the shrinkage that occurs when the polyester/styrene mixture crosslinks.\nThe development of low-profile unsaturated polyester compounds has led to a wide acceptance of these materials by the transportation industry because of their good surface appearance, dimensional stability, physical properties, assembly consolidation and potential weight savings. However, as new applications developed standards have been raised making it desirable for even better surface appearance and the elimination of ripples and waviness that sometimes develop, particularly in relatively large appearance sensitive areas.\nThere is, therefore, a need to provide low-profile unsaturated polyester compounds which afford improved surface appearance in the molded parts obtainable therefrom."}
{"text": "1. Field of the Invention\nThe present invention relates to computing. More particularly, the present invention relates to delayable events in a home network.\n2. Description of the Related Art\nUniversal Plug and Play (UPnP) is a distributed, open networking architecture that allows devices to connect seamlessly and to simplify the implementation of networks in the home (data sharing, communications, and entertainment) and corporate environments. UPnP achieves this by defining and publishing UPnP device control protocols built upon open, Internet-based communication standards.\nUPnP has grown in popularity of late in part due to the rise in popularity of media servers. Media servers are small computers that store multiple types of content (e.g., photos, music, videos, etc.). The content may then be streamed from a media server to one or more control points (e.g., iPod, television set, etc.).\nAs an example, a “Media Server” device might contain a significant portion of the homeowner's audio, video, and still-image library. In order for the homeowner to enjoy this content, the homeowner must be able to browse the objects stored on the Media Server, select a specific one, and cause it to be “played” on an appropriate rendering device.\nFor maximum convenience, it is highly desirable to allow the homeowner to initiate these operations from a variety of User Interface (UI) devices. In most cases, these UI devices will either be a UI built into the rendering device, or a stand-alone UI device such as a wireless PDA or tablet. In other cases, the home network user interface device could be more remote and communicate with the home network through a tunneling mechanism on the Internet.\nCurrently, UPnP clients (“control points”) may subscribe to event messages from a UPnP service. When an event occurs in the service, such as the addition or deletion of a file, service errors, etc., the subscribing control points receive event messages notifying them of the event. Each event reflects a change in some state variable being monitored by the service.\nThe drawback of this approach is that the service can only send events to control points that are connected to the network at the time the event is detected. The event publisher discards the event message if the subscriber doesn't respond for 30 seconds after the event transmission. If a control point is disconnected, or is in sleep mode, the control point never receives the event message."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a portable terminal, and more particularly, to a portable terminal provided with a wireless charging module.\n2. Description of the Related Art\nPortable terminals, such as a cellular phone and a smart phone, continue to grow in terms of functions as multimedia services are expanded. In addition, as various application programs are provided, user interface environments are developed for user convenience and to satisfy various user preferences.\nPortable terminals generally include bar-type, folder-type, sliding-type, and swing-type terminals. When mobile communications such as voice and short message transmissions were principal functions, folder-type and slider-type terminals were predominant in the market. However, enlarged display devices of portable terminals have increased as multimedia services have progressed. Accordingly, physical keypads have generally been replaced by a touch screen function, which improves the portability of a portable terminal since a display device may be enlarged and the portable terminal thickness may be reduced.\nIn addition, a portable terminal is provided with antenna devices that enable communication in various frequency bands such as a DMB (Digital Multimedia Broadcasting), a LAN (Local Area Network), an NFC (Near Field Communication), and a Bluetooth® antenna in addition to an antenna device for a wireless communication function. Recently, portable terminals have been equipped with an antenna and a module that provides a mobile charging function.\nFIG. 1 illustrates a portable terminal 100. In particular, FIG. 1 illustrates a configuration provided with a wireless charging module that includes a receiving-side resonant antenna 131 providing a wireless charging function and a receiving circuit unit provided on a substrate 133.\nAs illustrated in FIG. 1, the terminal 100 includes a battery-mounting groove 119 formed on a back surface of a body 101, and a camera module 117 provided at a side of the battery-mounting groove 119. The battery-mounting groove 119 is concealed by a cover member 102 detachably provided on the back surface of the body 101. A user may open the battery-mounting groove 119 by removing the cover member 102 as needed. In addition, the cover member 102 is provided with an opening 127 that exposes the camera module 117 so that a subject may be photographed even when the cover member 102 is coupled to the body 101.\nA wireless charging module is provided on the inner surface of the cover member 102. The wireless charging module includes a receiving-side resonant antenna 131 and a receiving circuit unit. A second cover member 141 may be provided on the inner surface of the cover member 102 so as to provide a stable installment structure of the receiving-side resonant antenna 131 and the receiving circuit unit.\nThe receiving-side resonant antenna 131 produces signal power by a magnetic induction or magnetic resonance phenomenon, according to an electromagnetic field generated at a primary coil of a charger (not illustrated), and transmits the signal power to the receiving circuit unit. A connection piece 131a connected to the receiving circuit unit is formed at a side of the receiving-side resonance antenna 131, which may be attached on the inner surface of the cover member 102. In the portable terminal 100 illustrated in FIG. 1, however, the second cover member 141 is coupled to the cover member 102 when the receiving-side resonant antenna 131 is attached to the inner surface of the second cover member 141. As a result, the receiving-side resonant antenna 131 is disposed on the inner surface of the cover member 102.\nThus, the terminal 100 is provided with an electromagnetic shielding member 139 in order to cut off an electromagnetic field effect exerted on circuit devices inside the body or a battery pack, formed around the receiving-side resonance antenna 131. The electromagnetic shielding member 139 is attached on the cover member 102 and interposed between the receiving-side resonant antenna 131 and the body 101. Since the receiving-side resonant antenna 131 is directly attached to the second cover member 141, the receiving-side resonant antenna 131 may be attached on the electromagnetic shielding member 139 after the electromagnetic shielding member 139 is attached to the second cover member 141 in advance.\nSince the receiving circuit unit includes the wireless charging circuit provided on the substrate 133, the receiving circuit unit converts signal power received through the receiving-side resonant antenna 131 into charging power, and provides the charging power to the battery pack mounted on the body 101. A connection portion 131b corresponding to the connection piece 131a is provided at a side of the substrate 133. The receiving circuit unit is also attached to the second cover member 141 together with the receiving-side resonant antenna 131. When the receiving-side resonant antenna 131 and the receiving circuit unit are attached to the second cover member 141, the connection piece 131a and the connection portion 131b are engaged with each other to be electrically connected.\nA flexible printed circuit board 135 is disposed at a side of the receiving circuit unit and includes a connection pad 137 at one end. The charging power provided from the receiving circuit unit is transmitted to the body 101 through the flexible printed circuit board 135 and the connection pad 137. The body 101 is provided with connection terminals 115 at a side of the battery-mounting groove 119. When the cover member 102 is coupled to the body 101, the connection pad 137 is connected with the connection terminals 115, thereby providing the charging power to the body 101, in particular, to the battery pack mounted on the body 101.\nFIG. 2 schematically illustrates a configuration of the wireless charging module disposed on the cover member 102. In FIG. 2, the height h of the wireless charging module, i.e. the thickness is illustrated. It is noted that the second cover member 141 is not illustrated in FIG. 2, and the wireless charging module is illustrated as being directly attachable to the cover member 102.\nThe receiving-side resonance antenna 131 and the electromagnetic shielding member 139 are stacked on the inner surface of the cover member 102. The receiving circuit unit includes circuit elements 133a such as a charging control circuit chip and an inductor, which are mounted on the substrate 133. The thicknesses of the receiving-side resonant antenna 131 and the substrate 133 are about 0.35 mm, the thickness of the electromagnetic shielding member 139 is about 0.6 mm, and the maximum height of the circuit elements 133a is about 1.25 mm. Accordingly, the maximum thickness of the wireless charging module from the inner surface of the cover member 102 is about 1.6 mm.\nAs described above, as multimedia functions of portable terminals have expanded, display devices have been enlarged and efforts have been made to reduce the thickness and weight of the portable terminals. A wireless charging module may enhance the convenience to charge a portable terminal but increases the thickness of the portable terminal. Particularly, slimmer portable terminals having a thickness of not more than 10 mm will have compromised size when equipped with a wireless charging module having a thickness of about 1.6 mm."}
{"text": "1. Field of the Invention\nThis invention relates to treadmills and, more particularly, the treadmills that have a tread base which is reorientable from a first exercise position to a second storage position within the cabinet, which cabinet includes latching structure for latching the tread base in the cabinet.\n2. State of the Art\nExercise treadmills typically include a frame having a left side and a right side spaced apart from the left side and in general alignment therewith. A rigid deck is also typically secured between the left side and the right side. A front roller and rear roller are typically connected to and extend between the left side and the right side forward and rearward of the deck. An endless belt is trained around the front roller and the rear roller. The user exercises on the treadmill by walking, jogging or running on the endless belt on top of a deck underlying the endless belt.\nTypical treadmills also include surface engaging structure to support the treadmill on a support surface. The surface engaging structure typically includes feet positioned proximate the rear of the treadmill and feet positioned proximate the front of treadmill. The front feet or the rear feet may be operable to vary the inclination of the treadmill with respect to the support surface. For example, U.S. Pat. No. 4,913,396 (Dalebout et al.) discloses a system for varying or adjusting the incline of a treadmill through the use of a pneumatic cylinder. U.S. Pat. No. 4,998,725 (Watterson et al.) discloses an alternate arrangement for varying the inclination of a treadmill.\nTreadmills also include handles or other upright structure such as that shown in U.S. Des. Pat. No. 304,849 (Watterson), U.S. Des. Pat. No. 306,468 (Watterson), U.S. Des. Pat. No. 306,891 (Watterson), U.S. Des. Pat. No. 316,124 (Dalebout et al.), U.S. Des. Pat. No. 318,699 (Jacobson et al.), U.S. Des. Pat. No. 323,198 (Dalebout et al.), and U.S. Des. Pat. No. 323,199 (Dalebout et al.). Reorientation or repositioning of the upright structure to facilitate storage has also been disclosed. U.S. Pat. No. 5,102,380 (Jacobson et al.) shows a treadmill in which a center post may be reoriented from an upright operating position to a lowered position in alignment with the treadmill and with the belt or deck. U.S. Des. Pat. No. 211,801 (Quinton) shows a treadmill with structure that may be moved from an upright position to a lowered position in general alignment with the treadmill belt or deck. U.S. Patent Des. 207,541 shows a treadmill that may be reoriented from a horizontal operating condition to an upright storage position.\nStoring exercise equipment inside a cabinet or other enclosure is also known. U.S. Pat. No. 4,300,761 (Howard) shows an exercise bench which may be repositioned interior a cabinet for purposes of storage. U.S. Pat. No. 3,741,538 (Lewis et al.) shows an arrangement in which the exercising structure is folded upright for storage against a wall surface. U.S. Pat. No. 3,642,279 (Cutter) shows a treadmill in which an upright structure may be reoriented to be generally in alignment with the endless belt for purposes of reorienting the treadmill to an upright or storage configuration.\nU.S. Pat. No. 4,679,787 (Guilbault) shows a bed combined with a treadmill or rolling structure in which the bed is positioned over the top of the treadmill or rolling structure for purposes of storage. U.S. Pat. No. 4,757,987 (Allemand) shows a treadmill which may be reconfigured into a compact foldable structure which may, in turn, be transported. U.S. Pat. No. 4,066,257 (Moller) shows a treadmill positioned within a cabinet that is secured to a wall and reoriented between an upright stored position and an extended or horizontal position for use."}
{"text": "Thin film transistors (TFTs) fabricated on polycrystalline silicon (polysilicon) have gained much attention in flat panel displays such as active matrix LCDs and in static random access memory units. In the future it can be expected that the degree of circuit integration will continue to increase as device characteristics improve further, so that an entire system will be formed on a single panel. In addition, in flat panel displays polysilicon thin film transistor technology enables the integration of row and column drive circuitry, and also additional functionality such as image reversal, aspect ratio control and level shifting among others. In addition to display elements and circuitry (both analog and digital), memories, solar cells, touch sensors and other sensors may all be integrated on the panel. For example, electrically erasable and programmable read only memories (EEPROMS) have been fabricated using a polysilicon TFT process.\nIn comparison with thick film devices TFT devices made with a thin film have the advantages of lower grain boundary trap density, higher mobility, and higher on-state current. It is desirable to make the thin film transistor as thin as possible in order to provide a high supply current in the on-state.\nHowever, thin film devices experience a higher lateral electric field, in particular at the junction between the channel and drain which arises from the reduced junction depth. This increase in lateral electric field is known to be a major cause of impact ionisation in the channel/drain region which results in the accumulation of holes. These holes are known to be the cause of a pronounced \"kink\" effect in the IV characteristics of a TFT device which in turn degrades the output characteristics of the device, and in particular reduces output resistance at high drain voltages and its gain, see for example M. Valdinoci. L. Colalongo, G. Baccarani. G. Fortunato, A. Pecora and I. Policicchio, \"Investigations on the Kink Effect in Poly-TFTs\", Proceediings of the ESSDERC, pp. 1055-1058, 1996 and A. G. Lewis, T. Y. Huang, R. H. Bruce, M. Koyangi, A. Chiang and I. W. Wu, \"Polysilicon Thin Film Transistor for Analogue Circuit Applications\", IEDM Tech. Digest, pp. 264-267, 1988. In addition the kink effect also causes avalanche induced short channel effects.\nFurthermore the high lateral electric field causes anomalous leakage current in the off state that counteracts the benefit provided by a thin film in which the lower trapped charge content would otherwise tend to reduce leakage current. This anomalous leakage current is a serious problem in poly-Si TFTs."}
{"text": "Field\nThe present disclosure relates to a wearable device, which provides haptic and audio feedback based on stereo camera input.\nDescription of the Related Art\nWearable cameras provide recording and documenting of a user's experience, often from the same or similar point of view or field of view (FOV) of the user. However, these devices are passive recorders, and do not provide real time processing and information about the scene in the FOV. Certain users, such as blind persons, may desire additional feedback relating to the environment. Other wearable cameras may be designed to assist blind persons. However, such devices lack stereo cameras for reliable depth perception information.\nThus, there is a need for an unobtrusive device which augments a user's environmental awareness with depth perception and object recognition."}
{"text": "In recent years, as disclosed in JP 5882522, for example, golf club heads have been proposed in which a raised portion is provided on the crown portion and a sloped surface is formed as a step between the raised portion and the portion rearward thereof. This configuration enables the height of the face portion to be raised by the height of the raised portion. Thus, the rebound performance of the face portion can be improved. Also, on the crown portion, only the raised portion is formed higher, and the portion rearward thereof is formed at a lower position than the raised portion, enabling the center of gravity of the head to be lowered.\nJP 5882522 is an example of related art."}
{"text": "This invention relates to a wiper device equipped with a rise-up mechanism.\nIn general, a rise-up mechanism is a mechanism for shifting a wiper blade to a position away from its normal wiping positions when it stops so that a driver does not see the blade in front of his eyes.\nFIGS. 1a to 1c show schematically a sequence of wiping operational steps. During the starting step, the blade 18 shifts from a parked position A through the position B to the position C. Thereafter, the blade 18 reciprocates between the positions B and C during its normal wiping operation. At the stopping step, the blade 18 shifts from the position B to the position C and then again from the position C through the position B to the parked position A in which it stops.\nFIGS. 2 to 4 show a conventional rise-up mechanism for a wiper device in which the effective length of a connecting rod or link is changed by reverse rotation of a motor when a wiper switch is switched off.\nA cover plate 3, latch spring 4, latch 5, motor arm 6 and eccentric member 7 are assembled in order onto a motor shaft 2 of a motor 1 and fixed thereto by means of a screw 10. The cover plate 3 is attached through a grommet 11 to the motor body. The eccentric member 7 is biased into the central axial hole of the motor arm 6 by means of a spring 12. In one end thereof, a stopper 8 is attached to engage the groove 7a of the eccentric member 7 and a spring 13 is attached to bias the stopper. At the other end thereof, a joining shaft 14 is provided to be joined to a connecting rod 15 is provided. A groove 6a is formed in the lower surface of the motor arm 6 near the joint shaft 14 and engages a pawl 5a of the latch 5. A pawl 7b projects from the lower surface of the eccentric member 7 and presses the inner actuating portion 5b of the latch 5 to slide the latch 5 so that the pawl 5a becomes disengaged from the groove 6a.\nIn operation, first, the motor 1 rotates to the right during starting. When the motor shaft 2 starts to rotate to the right, the pawl 5a of the latch 5, which is prevented from rotating by the cover plate 3, engages the groove 6a of the motor arm 6 as shown in FIG. 4, so that the motor 1 does not rotate while the motor shaft 2 and the eccentric member 7 rotate together as shown in FIG. 4b. When they rotate by 180.degree., as shown in FIG. 4c, the stopper 8 engages the groove 7a of the eccentric member 7 by the biasing force of the spring 13. At the same time, the pawl 76 of the eccentric metal 7 presses the inner actuating portion 5b of the latch 5 to slidably shift the latch 5 to the left as shown in FIG. 4c. As a result, the pawl 5a of the latch 5 leaves the groove 6a of the motor arm 6.\nAssuming that the length or distance between the center of the joint shaft 14 and the center of the eccentric member 7 is R and that the length (eccentric degree) between the center of the eccentric member 7 and the rotation center thereof is r, the effective length between the centers of the motor shaft 2 and the joint shaft 14 of the motor arm 6 changes from (R+r) to (R-r) due to the eccentric effect when the eccentric member 7 rotates by 180.degree..\nFrom the condition of FIG. 4c, the motor arm 6 rotates as shown in FIGS. 4d to 4f in response to the rotation of the motor shaft 2 together with the eccentric member 7 resulting from engaging stopper 8 with the groove 7a. The motor arm 6 further rotates and returns from the condition of FIG. 4f to the condition of FIG. 4c. While the motor arm 6 rotates as shown in FIGS. 4c to 4f, the connecting rod 15 joined to the joint shaft 14 reciprocates to swing the wiper pivot 16 around the axis 17 so that the blade 18 reciprocates in its normal wiping range between the positions B and C in FIG. 1b.\nWhen a wiper switch (not shown) is switched off to stop the motor shaft 2, it continues to rotate until the condition of FIG. 4c where the motor shaft 2 begins to rotate in a reverse direction. After such reverse rotation of the motor shaft 2 in the condition of FIG. 4c, the groove 6a of the motor arm 6 contacts the pawl 5a of the latch 5. (It does not contact the pawl 5a in case of the right rotation thereof). As shown in FIG. 4g, the groove 7a of the eccentric member 7 pushes the stopper 8 at its tapered surface during its reverse rotation against the biasing force of the spring 13 so that the stopper 8 is disengaged therefrom. Thus, as shown in FIG. 4g, the motor arm 6 does not rotate and only the eccentric member 7 rotates so that after its 180.degree. rotation the pawl 5a engages the groove 6a as shown in FIG. 4a. Also, the length or distance between the center of the joint shaft 14 of the motor arm 6 and the motor shaft 4 changes from (R-r) to (R+r) due to the eccentric effect of the eccentric member 7 by its 180.degree. rotation. The effective length of the connecting rod 15 changes so that the blade 18 can be received in the position A of FIG. 1.\nIn the above-stated conventional wiper device, however, the latch 5 slides during the rise-up operational step. The rise-up operation cannot be accomplished by rotation of the members only. For this reason, the operation is not sure and stable. In addition, as the rotating members have some sliding portions, the mechanism is apt to be complicated. Thus, assembling thereof is difficult and its production cost is high.\nOn the other hand, during the starting step, the blade 18 starts to operate only after the motor shaft 2 freely rotates by 180.degree.. Therefore, the response of the device is poor.\nAlthough not shown, in another conventional rise up mechanism for a wiper device, a motor itself moves to shift a motor shaft. In such a device, a mechanism for shifting the motor is complex in construction. As compared with a general wiper device, a rise up mechanism is expensive."}
{"text": "1. Field of Invention\nThe embodiment of the present invention relates generally to a transmitting method and, more particularly, to a method for transmitting data stream.\n2. Description of Related Art\nIn traditional package transmitting mechanism, when a transmitting end receives an ACK package transmitted from a receiving end, the transmitting end continuously transmits next package. That is to say, the transmitting end stops transmitting packages when the transmitting end does not receive the ACK package transmitted from the receiving end, or the transmitting end disconnects a communication with the receiving end directly when the transmitting end does not receive the ACK package transmitted from the receiving end for a period.\nIn addition, the bandwidth and the buffering space for transmitting the package are different owing to differences of a quality of a content of a film, the way to compress the film, and so on. When a user selects one of films, a client end download related streams of the film from a server. However, in this mode, there will be a bandwidth utilizing shake phenomenon happened in the server and the client end, and the quality of the film will be affected if the package disappears.\nMany efforts have been devoted trying to find a solution of the aforementioned problems. Nonetheless, there still a need to improve the existing apparatus and techniques in the art."}
{"text": "In the automotive and paper industries, the use of steel belts and chains is common for conveying heavy loads. These belts typically require lubrication which involves regular maintenance. Also, steel conveyors are usually not very smooth and may damage the conveyed products. It is therefore of interest to use heavy plastic modular belts or chains which do not require lubrication and are smooth and less hard on their conveying surface. Since modular belts and chains are supported by slider beds made from plastics or steel, the friction between the belt and the slider bed increases the required driving power and at the same time reduces the maximum conveyable load.\nAnother problem with heavy conveyors is the large belt thickness and sprocket diameter which typically requires a heavy and space-consuming conveyor construction. Often the conveying surface for such conveyors needs to be near the ground. This arrangement is specifically necessary for people mover belts. The purpose of such belts is to transport workers slowly along an assembly line on which the components to assemble are also moving. An application for this type of belt is found in the automobile manufacturing industry. Because the heavy people mover conveyors require space, they are required to be installed in the ground. In order to install the conveyor at or near ground level, a conveyor with a large height requires a deep pit. The same problems also arise in connection with the transportation of heavy goods such as paper rolls.\nAccordingly, there is a need for a transport system that reduces the friction forces to improve the loading capacity and reduce the depth of the pit in the ground that may be needed in applications where the conveyor is installed at or near ground level."}
{"text": "The performance of silicon-based complementary metal-oxide-semiconductor (CMOS) transistors steadily improves as device dimensions shrink. The decreasing size of metal-oxide-semiconductor field-effect transistors (MOSFETs) provides improved integrated-circuit performance speed and cost per function. As channel lengths of MOSFET devices are reduced to increase both the operation speed and the number of components per chip, the source and drain regions extend towards each other, occupying the entire channel area between the source and the drain. Interactions between the source and drain of the MOSFET degrade the ability of the gate of the MOSFET to control whether the MOSFET is “on” or “off.” In particular, the threshold voltage and drive current decrease appreciably with the channel length. This phenomenon is called the “short channel effect.” The term “short channel effect,” as used herein, refers to the limitations on electron drift characteristics and modification of the threshold voltage caused by shortening trench lengths.\nDouble- or tri-gate transistors, such as vertical double-gate silicon-on-insulator (SOI) transistors or fin-FETs, offer significant advantages related to high drive current and high immunity to short trench effects. Conventionally, fin-FET devices have included single, unitary semiconductor structures that protrude from an active surface of a substrate. Such a semiconductor structure is generally referred to as a “fin.” A polysilicon layer may be deposited over a central portion of the fin and patterned to form a pair of gates on opposite sides of the fin. Among the many advantages offered by fin-FETs is better gate control at short gate lengths. Fin-FETs facilitate down-scaling of CMOS dimensions while maintaining acceptable performance.\nWith ever-decreasing semiconductor device feature sizes, the effects of shortened channel lengths become increasingly problematic in the fabrication of semiconductor devices.\nMethods of fabricating semiconductor devices to reduce short channel effects and increase drive current, as well as improved fin-FET structures, are desirable."}
{"text": "1. Field of the Invention\nThis invention relates to solar energy assemblies which are adapted to provide decorative and functional means for collecting radiant solar energy and, more specifically, it is directed toward unique panel constructions adapted for such purposes.\n2. Description of the Prior Art\nVarious forms of functional and decorative building construction components positioned on building exteriors such as vertical exterior walls and roofs have been known for years. Not only has it been known to provide decorative wall coverings for the interior, but various forms of exterior siding have been known. See, for example, U.S. Pat. Nos. 2,642,968, 2,777,549, 3,054,223 and 3,394,520.\nAs a result of the shortage of energy on a worldwide basis, more and more effort is being directed toward more efficient use of existing energy supplies. For example, in order to conserve our coal, gas and oil reserves more emphasis has been placed upon maintaining of residential and commercial structures at reduced temperatures in cold weather and providing increased thermal insulation to minimize heat loss. There has also been a great deal of emphasis directed toward the use of solar energy in heating of buildings, heating of hot water and other uses.\nU.S. Pat. No. 3,918,430 discloses a hot water solar system adapted for use on a roof or other portion of a building. A plurality of water channels are housed within a rigid frame underlying a series of layers of plastic material.\nU.S. Pat. No. 4,029,080 discloses a thermal collector for a solar energy system. The prime thrust of this disclosure is directed toward an air system adapted for use on a roof.\nU.S. Pat. No. 4,069,809 discloses a solar system wherein a series of building blocks have transparent members for permitting passage of the sun's rays therethrough. The series of blocks provides a vertical air channel passing immediately behind the transparent window in each block and a series of three generally vertically oriented passageways positioned within each block remote from the front transparent window.\nU.S. Pat. No. 4,120,282 discloses a solar system consisting of a number of fixed flat plate solar reflectors and collectors.\nU.S. Pat. No. 4,073,282 discloses a solar collecting system wherein a matrix of expanded sheets having large openings is employed to collect the sun's radiant energy. Means are provided for circulating air through the chamber and into contact with the slit and expanded sheets.\nU.S. Pat. Nos. 4,076,015 and 4,077,393 each disclose systems wherein modular elements provide a plurality of raised surfaces for receipt of the sun's rays as used in combination with raised reflective surfaces. Among the problems encountered with known solar collecting systems are the somewhat unsightly nature of the same and, in some instances, the expense of installing the same.\nThere remains a need for a solar collecting system for exterior walls, roofs and other portions of buildings which is both decorative and functional. There is a further need for such systems which can be applied readily to preexisting buildings as well as buildings designed and constructed with the solar energy system in mind."}
{"text": "1. Field of the Invention\nThe invention relates to an eyelet arrangement for use with a swimming pool having a plastic liner which permits the part having the eye to be removed to allow changing of the plastic liner without requiring access to the outside of the wall.\n2. The Prior Art\nThe patent to Engelhart U.S. Pat. No. 3,868,732 shows a device for accomplishing the same function as the present invention. The device of Engelhart however is quite expensive, as it involves the use of an externally and internally threaded sleeve having a large flange at one end, the making of which requires a number of costly operations."}
{"text": "It is known that the content of unsaturated fatty acid in beef fat is associated with quality of beef, such as the beef taste and texture. Generally, if the content of unsaturated fatty acid is high with low melting point, the beef is considered to have good quality, which gives good taste and texture.\nHowever, in order to judge whether the beef has such quality, there was no other way but to depend on the subjective (sensuous) method, in which the beef was actually eaten and evaluated. In other words, there was no conventional method, which was more objective, simple and efficient, like the judgment approach of predicting quality of beef, simply based on the genotype of a specific bovine gene.\nBy the way, stearoyl-CoA desaturase (SCD) is known as an enzyme which desaturates beef fat. This enzyme desaturates stearoyl-CoA which plays an important role for in vivo biosynthesis of lipids and their degradation. Amino acid sequence and cDNA sequence of bovine stearoyl-CoA (derived from Bos taurus) are shown in DDBJ/EMBL/GenBank databases; Accession number “AB075020”. Information of these sequences was provided by inventors of the present invention.\nAs mentioned above, the content of unsaturated fatty acid in beef fat (in other words, the amount of the unsaturated fatty acid content) is associated with quality of beef, such as the beef taste and texture. If the amount of the unsaturated fatty acid content can be predicted on the basis of the genotype of a specific bovine gene, then it provides a new method which enables simple examination as to quality of beef, such as the beef taste and texture.\nSuch method based on the genotype is useful not only for evaluation and selection of cattle (beef cattle) with good quality of beef, but also for breeding and reproduction of cattle.\nFurthermore, the content of unsaturated fatty acid in beef fat is associated with the cholesterol accumulation caused by beef intake, and it is considered one of important features of the cattle, especially in the countries where beef is eaten very often, such as the Europe and the United States of America. Therefore, the above-mentioned method for evaluating the unsaturated fatty acid content based on the genotype, is also useful from the view point of health care.\nIn addition, the content of unsaturated fatty acid in milk fat of the dairy cattle is considered to affect the taste and mouthfeel of dairy products, such as butter. For example, when the content of unsaturated fatty acid in milk fat is higher and the melting point is lower, then the butter produced by using the milk as raw material is softer with good mouth-melt. Therefore, the above-mentioned method for predicting the unsaturated fatty acid content based on the genotype, is also useful for dairy products and dairy cattle breeding."}
{"text": "This invention pertains to the field of digital computing and, more particularly, to a large scale digital computing system capable of performing high-speed addition without unduly increasing power consumption.\nAddition is a fundamental operation in the arithmetic unit of a digital computing system because it is employed to perform, not only the addition function, but also the multiplication function and the accumulation function. In a conventional adder, the processing time is directly related to the number of bits in the numbers being processed because of the necessity for carries to propagate from stage to stage before the arithmetic result is generated. A number of schemes have been developed over the years, such as the so-called look ahead adder, to generate arithmetic results without the delays occasioned by carry propagation. These schemes require a great deal of extra electronic circuitry and, therefore, dissipate more power than a conventional adder."}
{"text": "1. Field of the Invention\nThe present invention relates to semiconductor products reliability testing and, more particularly, to semiconductor reliability test chips for testing standard and ASIC semiconductor packages.\n2. Background and Related Art\nDuring the course of qualifying packages and modules, it is customary to run standard stresses to predict reliability of the packaged semiconductor products under field conditions. Typically, the reliability of the semiconductor packages is tested by subjecting them to a variety of life accelerating environments over a period of time until product failure or minimum requirements are met. The packages are then inspected and tested in an attempt to determine the cause of failure. Since there can be many reasons for failure, the analyses of the failure can be lengthy and difficult. Attempts have also been made to design semiconductor test chips to assess specific types of failure of the product.\nFor example, an article by J. S. Sweet, entitled “The Use of Special Purpose Assembly Test Chips for Evaluating Reliability In Packaged Devices”, published by Sandie National Laboratory, pages 15–19, describes some of these types of chips. The article describes a series of individual special purpose assembly test chips to aid in assessing the reliability of packaged integrated circuits. The special purpose assembly test chips contain special purpose circuits or sensors which enhance the detection of failures or detect moisture, detect mobile ions, or other contaminants which can lead to failure of the semiconductor component.\nOther special purpose test chips have been designed to aid in assessing the reliability of a variety of specific types of failures of semiconductor packages. For example the U.S. Pat. No. 6,538,264 to Corbett, et al. describes a test chip with a plurality of test functions, such as, bond pad pitch and size effects on chip design, wire bond placement accuracy, bond pad damage below the bond pad during bonding (cratering), street width effects, thermal impedance effects, ion mobility evaluation and chip on board in flip chip application test capabilities.\nTest chips for flip-chip packages, such as described by Corbett, et al., using C4 solder ball technology have thus taken a variety of forms. The need to verify C4 integrity as part of the chip/package/interconnect qualification is an important product requirement. This requirement has become more important with the advent of organic C4 chip packages. In this regard, it has been found that certain product design features will result in early reliability stressing failures. Thus, to provide the most effective testing process, it is necessary to design the chip/package/interconnect qualification packages as closely to product as possible to avoid having to address, either failures in features that do not appear in the product, or failure to stress features that do exist in the product that may ultimately fail. One of the major stresses that cause failures in organic flip chip packages is the thermal mismatch in CTE between the flip chips and organic substrate. Such thermal mismatch causes stress and, potentially, fatigue at the C4 interconnect initiating fracture and cracking of the C4 bonding, for example, resulting in connection failure. In this regard, it is known that this stress is proportional to the distance from neutral point (DNP) of a particular C4 solder ball connection.\nOne approach to stressing packages to test for fatigue leading to fracture and cracking of C4 connections due to thermal mismatch is to cyclically heat the packaged chip using electrically resistive heaters in the chip to simulate product thermal cycling. This can be accomplished by designing a test chip with a large resistive heater in the chip M1 metal layer. Such a heater is typically wired through a small number of C4 connections. With a small number of connections and with the need for increased heat and power, there is concern with the resistive heating of the relatively narrow package signal traces such as to potentially introduce excessive temperature induced failures that would not exist in the product. Alternatively, designing the package wires as heavy power supplies connected to signal C4 positions represents a nonstandard feature in the package that does not represent product.\nA conventional approach to testing for the reliability and integrity of C4 interconnections between chip and substrate is to employ a continuity-type testing procedure. An example of such an approach is that employed by Corbett, et al. supra wherein metal wire stitch lines are employed to connect selected C4 pads on the chip together and substrate or board level wiring is used to connect all wiring in a daisy chain approach. Thus, the resulting structure has a concatenation of board wire, package wire, chip wire, board wire, etc. While this approach has the advantage of allowing a large number of connections to be tested with one circuit, it has the disadvantage that should there be fatigue or fracture in one of the interconnections being monitored causing a change in resistance, for example, such change can easily be lost in the larger overall resistance of the single circuit interconnecting all of the interconnects.\nIn this regard, it is known that small changes in resistance are indicative of C4 fatigue and crack initiation. Accordingly, it is advantageous to test for fatigue and crack initiation using a low resistance circuit approach such as to allow easy and ready identification of the connection exhibiting fatigue and crack initiation.\nA further limitation of prior art approaches is that the array of pads on the test chip used for testing covers a small area of the chip, and the array of pads is typically near the center of the chip thus discounting the contribution to stress that would be expected for high DNP C4's.\nAccordingly, the test chip should be made to replicate the product that it is representing as closely as possible. To this end, it is undesirable to wire out all of the C4 pads to the chip substrate or PCB. In this regard, most product chips require power connections that are handled via the power planes of the package. Thus, to maintain the mechanical properties of the package, it is advantageous to design the test package such that signal and power structures look like the product design. It should be noted that, the highest DNP C4's are typically power and ground connections, particularly in application specific footprints."}
{"text": "It has long been known that absorbent articles such as conventional absorbent articles (e.g., diapers, adult incontinence articles, feminine hygiene pads) offer the benefit of receiving and containing urine and/or other bodily exudates (e.g., feces, menses, mixture of feces and urine, mixture of menses and urine, etc.). To effectively contain bodily exudates, the article should provide a snug fit around the waist and legs of a wearer.\nCurrent diaper designs frequently include the use of a barrier leg cuff to prevent leakage of bodily exudates and an outer cuff which provides a covering over the barrier leg cuff to minimize the visibility of exudates through the barrier cuff and provide a secondary means to capture bodily exudates should they breach the barrier leg cuff. The barrier leg cuff may be made using a hydrophobic nonwoven and may be disposed on the body-facing surface of the absorbent article or connected to the body-facing surface of the film backsheet layer. The barrier leg cuff may be a substantially liquid impervious layer that prevents bodily exudates from passing out of the sides of the article and may also be highly breathable, allowing outside air to reach the skin to help maintain a healthy level of skin hydration. In many current diapers, the outer cuff comprises the polymeric film layer of the backsheet to provide high opacity required to cover the barrier leg cuff as well as to prevent molten adhesive from passing through the cuff to the garment-facing surface of the article during manufacturing. The outer cuff contains the outer leg elastic strands, which create the contraction forces and gathers, and can be sandwiched between the cuff material and backsheet material. The elastic strands in the leg cuffs are typically joined with molten adhesive during manufacture, and the hot adhesive generally has the potential to pass through nonwoven materials during manufacture, causing contamination of manufacturing lines as well as the potential for stickiness on the outside surface of the article. The polymeric film generally is used to prevent these issues, however, results in a plastic-like look as well as a noisy application process.\nBecause of manufacturing tolerances when cutting, tracking, and combining materials, the outer leg elastic strands are generally spaced inboard from the longitudinal edge of the article in the crotch region. This prevents inadvertent cutting or exposure of the outer leg elastic strands during the manufacturing process. This design does not result in the outermost portion of the longitudinal edge of the product continuously contacting closely to the skin of the user during wear. Thus, the ability of the elastic strand(s) to control the edge of the article diminishes as the distance between the outermost elastic and the edge increases, leading to a more random distribution of larger gathers which contact the skin at larger intervals or sometimes not at all. This effect can lead to user perception that the diaper may leak where the longitudinal edge does not contact the skin of the user. In addition, many articles currently available contain only two to three outer leg elastics per side to create the gathers, increasing the difficulty of achieving the desired appearance of a wide finished leg cuff or more garment-like cuff such as the elasticized hemmed edge of the arm cuff of a sweater. If the elastics are spaced more closely, the result is a narrow section of elasticized zone, which results in a less finished, less comfortable, and less clothing-like appearance. If the elastics are spaced farther apart, the gathers can appear to separate further from the skin of the user, leading to a perception of potential leakage risk. As discussed above, this is driven by having less control of the gathers between strands of increasing separation.\nAccordingly, it is desirable to provide an absorbent article with a folded outer cuff design having finished edges with elastics that are close to the edge to maintain a close proximity to the skin to create improved fit, a more aesthetically pleasing, clothing-like design and improved leakage protection.\nHowever, even with the improved leakage protection provided by the cuff designs detailed herein, the most common mode of failure for absorbent articles still occurs when body exudates leak out of the gaps between the article and the wearer's legs and/or waist. When fecal material (e.g., runny bowel movement, a mixture of bowel movement and urine, etc.) is not absorbed into the topsheet and core of absorbent article, the fecal material can leak out of the gaps between the article and the wearer's legs or waist. In situations where a wearer exudes a higher quantity of fecal material—which is absorbed by the absorbent core more slowly than urine—the fecal material may move laterally along the body-side surface of the absorbent article and reach the barrier leg cuff. After the fecal material reaches the barrier leg cuffs, it may travel longitudinally along the barrier leg cuffs. Due to the movement of the wearer and/or a shortage of available space under the barrier cuffs and/or within the absorbent article, the fecal material may leak out of the gaps between the article and the wearer's legs and/or waist. This results in soiling, wetting, or otherwise contaminating the wearer's clothing or other articles (e.g., bedding, furniture, caregiver clothing, etc.) that come in contact with the wearer's leaky absorbent article.\nAccordingly, it is of continued interest to provide an economically viable disposable absorbent article with the ability to minimize the negative effects of bodily extrudate leaks, while also making it easier to clean the wearer when the soiled disposable absorbent article is removed. To that end, it is of continued interest to provide a disposable absorbent article having sufficient retention capability to safely and cleanly retain bodily extrudate away from the wearer's clothing and/or skin throughout the expected time of article use."}
{"text": "I. FIELD OF THE INVENTION\nThe invention relates generally to a metal discriminator, and specifically to a discriminating circuit including a sense coil and means for measuring a change in the self inductance of the coil caused by a metallic object placed substantially in the center of the coil.\nII. RELATED ART\nPrevious discriminator systems based on the principle of applying pulsed signals into a sense coil and measuring the loss in voltage potential across the coil versus time that occurs in the coil after the pulse is applied are known. However, these systems tend to be very sensitive to external noise sources and are subject to temperature sensitivity and physical parameters such as cable length."}
{"text": "Prions are infectious pathogens that cause central nervous system spongiform encephalopathies in humans and animals. Prions are distinct from bacteria, viruses and viroids. The predominant hypothesis at present is that no nucleic acid component is necessary for infectivity of prion protein. Further, a prion which infects one species of animal (e.g., a human) will not infect another (e.g., a mouse).\nA major step in the study of prions and the diseases that they cause was the discovery and purification of a protein designated prion protein (\"PrP\") [Bolton et al., Science 218:1309-11 (1982); Prusiner et al., Biochemistry 21:6942-50 (1982); McKinley et al., Cell 35:57-62 (1983)]. Complete prion protein-encoding genes have since been cloned, sequenced and expressed in transgenic animals. PrP.sup.C is encoded by a single-copy host gene [Basler et al., Cell 46:417-28 (1986)] and is normally found at the outer surface of neurons. A leading hypothesis is that prion diseases result from conversion of PrP.sup.C into a modified form called PrP.sup.Sc.\nIt appears that the scrapie isoform of the prion protein (PrP.sup.Sc) is necessary for both the transmission and pathogenesis of the transmissible neurodegenerative diseases of animals and humans. See Prusiner, S. B., \"Molecular biology of prion disease,\" Science 252:1515-1522 (1991). The most common prion diseases of animals are scrapie of sheep and goats and bovine spongiform encephalopathy (BSE) of cattle [Wilesmith, J. and Wells, Microbiol. Immunol. 172:21-38 (1991)]. Four prion diseases of humans have been identified: (1) kuru, (2) Creutzfeldt-Jakob Disease (CJD), (3) Gerstmann-Strassler-Scheinker Disease (GSS), and (4) fatal familial insomnia (FFI) [Gajdusek, D.C., Science 197:943-960 (1977); Medori et al., N. Engl. J. Med. 326:444-449 (1992)]. The presentation of human prion diseases as sporadic, genetic and infectious illnesses initially posed a conundrum which has been explained by the cellular genetic origin of PrP.\nMost CJD cases are sporadic, but about 10-15% are inherited as autosomal dominant disorders that are caused by mutations in the human PrP gene [Hsiao et al., Neurology 40:1820-1827 (1990); Goldfarb et al., Science 258:806-808 (1992); Kitamoto et al., Proc. R. Soc. Lond. 343:391-398. Iatrogenic CJD has been caused by human growth hormone derived from cadaveric pituitaries as well as dura mater grafts [Brown et al., Lancet 340:24-27 (1992)]. Despite numerous attempts to link CJD to an infectious source such as the consumption of scrapie infected sheep meat, none has been identified to date [Harries-Jones et al., J. Neurol. Neurosurg. Psychiatry 51:1113-1119 (1988)] except in cases of iatrogenically induced disease. On the other hand, kuru, which for many decades devastated the Fore and neighboring tribes of the New Guinea highlands, is believed to have been spread by infection during ritualistic cannibalism [Alpers, M. P., Slow Transmissible Diseases of the Nervous System, Vol. 1, S. B. Prusiner and W. J. Hadlow, eds. (New York: Academic Press), pp. 66-90 (1979)].\nThe initial transmission of CJD to experimental primates has a rich history beginning with William Hadlow's recognition of the similarity between kuru and scrapie. In 1959, Hadlow suggested that extracts prepared from patients dying of kuru be inoculated into nonhuman primates and that the animals be observed for disease that was predicted to occur after a prolonged incubation period [Hadlow, W. J., Lancet 2:289-290 (1959)]. Seven years later, Gajdusek, Gibbs and Alpers demonstrated the transmissibility of kuru to chimpanzees after incubation periods ranging form 18 to 21 months [Gajdusek et al., Nature 209:794-796 (1966)]. The similarity of the neuropathology of kuru with that of CJD [Klatzo et al., Lab Invest. 8:799-847 (1959)] prompted similar experiments with chimpanzees and transmissions of disease were reported in 1968 [Gibbs, Jr. et al., Science 161:388-389 (1968)]. Over the last 25 years, about 300 cases of CJD, kuru and GSS have been transmitted to a variety of apes and monkeys.\nThe expense, scarcity and often perceived inhumanity of such experiments have restricted this work and thus limited the accumulation of knowledge. While the most reliable transmission data has been said to emanate from studies using nonhuman primates, some cases of human prion disease have been transmitted to rodents but apparently with less regularity [Gibbs, Jr. et al., Slow Transmissible Diseases of the Nervous System, Vol. 2, S. B. Prusiner and W. J. Hadlow, eds. (New York: Academic Press), pp. 87-110 (1979); Tateishi et al., Prion Diseases of Humans and Animals, Prusiner et al., eds. (London: Ellis Horwood), pp. 129-134 (1992)].\nThe infrequent transmission of human prion disease to rodents has been cited as an example of the \"species barrier\" first described by Pattison in his studies of passaging the scrapie agent between sheep and rodents [Pattison, I. H., NINDB Monograph 2, D. C. Gajdusek, C. J. Gibbs Jr. and M. P. Alpers, eds. (Washington, D.C.: U.S. Government Printing), pp. 249-257 (1965)]. In those investigations, the initial passage of prions from one species to another was associated with a prolonged incubation time with only a few animals developing illness. Subsequent passage in the same species was characterized by all the animals becoming ill after greatly shortened incubation times.\nThe molecular basis for the species barrier between Syrian hamster (SHa) and mouse was shown to reside in the sequence of the PrP gene using transgenic (Tg) mice [Scott et al., Cell 59:847-857 (1989)]. SHaPrP differs from MoPrP at 16 positions out of 254 amino acid residues [Basler et al., Cell 46:417-428 (1986); Locht et al., Proc. Natl. Acad. Sci. USA 83:6372-6376 (1986)]. Tg(SHaPrP) mice expressing SHaPrP had abbreviated incubation times when inoculated with SHa prions. When similar studies were performed with mice expressing the human, or ovine PrP transgenes, the species barrier was not abrogated, i.e., the percentage of animals which became infected were unacceptably low and the incubation times were unacceptably long. Thus, it has not been possible, for example in the case of human prions, to use transgenic animals (such as mice containing a PrP gene of another species) to reliably test a sample to determine if that sample is infected with prions. The seriousness of the health risk resulting from the lack of such a test is exemplified below.\nMore than 45 young adults previously treated with HGH derived from human pituitaries have developed CJD [Koch et al., N. Engl. J. Med. 313:731-733 (1985); Brown et al., Lancet 340:24-27 (1992); Fradkin et al., JAMA 265:880-884 (1991); Buchanan et al., Br. Med. J. 302:824-828 (1991)]. Fortunately, recombinant HGH is now used, although the seemingly remote possibility has been raised that increased expression of wtPrP.sup.C stimulated by high HGH might induce prion disease [Lasmezas et al., Biochem. Biophys. Res. Commun. 196:1163-1169 (1993)]. That the HGH prepared from pituitaries was contaminated with prions is supported by the transmission of prion disease to a monkey 66 months after inoculation with a suspect lot of HGH [Gibbs, Jr. et al., N. Engl. J. Med. 328:358-359 (1993)]. The long incubation times associated with prion diseases will not reveal the full extent of iatrogenic CJD for decades in thousands of people treated with HGH worldwide. Iatrogenic CJD also appears to have developed in four infertile women treated with contaminated human pituitary-derived gonadotrophin hormone [Healy et al., Br. J. Med. 307:517-518 (1993); Cochius et al., Aust. N.Z. J. Med. 20:592-593 (1990); Cochius et al., J. Neurol. Neurosurg. Psychiatry 55:1094-1095 (1992)] as well as at least 11 patients receiving dura mater grafts [Nisbet et al., J. Am. Med. Assoc. 261:1118 (1989); Thadani et al., J. Neurosurg. 69:766-769 (1988); Willison et al., J. Neurosurg. Psychiatric 54:940 (1991); Brown et al., Lancet 340:24-27 (1992)]. These cases of iatrogenic CJD underscore the need for screening pharmaceuticals that might possibly be contaminated with prions.\nRecently, two doctors in France were charged with involuntary manslaughter of a child who had been treated with growth hormones extracted from corpses. The child developed Creutzfeldt-Jakob Disease. (See New Scientist, Jul. 31, 1993, page 4). According to the Pasteur Institute, since 1989 there have been 24 reported cases of CJD in young people who were treated with human growth hormone between 1983 and mid-1985. Fifteen of these children have died. It now appears as though hundreds of children in France have been treated with growth hormone extracted from dead bodies at the risk of developing CJD (see New Scientist, Nov. 20, 1993, page 10.) In view of such, there clearly is a need for a convenient, cost-effective means for removing prions which cause CJD from blood and blood products. The present invention provides such a method."}
{"text": "Fence installers often set a first post in cement, insert a panel into the first post, insert the panel into a second post, set the second post in cement, then go back and screw the panel to the first post. As a general practice, these tasks are completed while the cement is still wet. In the frost-belt, this type of installation requires a longer post. By installing this way, there is always a chance that someone will lean on or bump into the panels, moving the posts and causing the posts to be crooked when the cement hardens. The posts take multiple steps to manufacture. Posts must be loaded into a machine twice; once to punch/route rail holes and once to drill screw holes. Hundreds of screws, longer posts, additional manufacturing steps, and increased installation time all increase the cost of the fence.\nSometimes panels need to be made shorter while on a job site, making it necessary to notch the end of the rails so that the rails will fit properly into their respective posts. If a panel needs to be shortened on the job site, the installer has a choice to either purchase a special tool (costing $200.00-$300.00) or try to notch the rails using a hack saw. Using a hack saw to notch the rails is very difficult, and increases costs and time on the job.\nThe panels of fences assembled using the conventional technique are not easily removable. If a panel is damaged, the posts must be removed to perform repairs. Consequently, repairs are extremely expensive. Another problem for consumers is insects (e.g., bees, etc.) building hives in the fence posts, because most aluminum fence products have no system in place to keep the insects out.\nIt would be desirable to implement a clip facilitating quick assembly and disassembly of fence components."}
{"text": "The present application claims priority of Japanese patent application Ser. No. 82/156399, filed Sept. 8, 1982.\nThe present invention relates to novel foam-forming, room temperature vulcanizable (RTV) organopolysiloxane compositions having increased foam-forming ratio and excellent foamability. Some technologies have been known for RTV foam-forming organopolysiloxane compositions, but most of them could not be applied to commercial production.\nIn U.S. Pat. No. 3,070,555 Bruner invented RTV foam-forming organopolysiloxane compositions utilizing the dehydrogenation reaction of .alpha.,.omega.-diorganopolysiloxane diol and organohydrogenpolysiloxane in the presence of a stannous carboxylate catalyst.\nIn U.S. Pat. No. 3,338,847 Nitzsche invented similar compositions comprising .alpha.,.omega.-diorganopolysiloxane diol, organohydrogenpolysiloxane, and unsaturated hydrocarbon containing hydroxyorganopolysiloxane in the presence of a metal carboxylate catalyst.\nIn U.S. Pat. No. 2,956,032 Joyce invented flame retardant RTV foam-forming organopolysiloxane compositions comprising .alpha.,.omega.-diorganopolysiloxane diol, organohydrogenpolysiloxane, stannous carboxylate and NiBr.sub.2.\nIn U.S. Pat. No. 3,428,580 Nitzsche invented flame retardant RTV foam-forming organopolysiloxane compositions comprising .alpha.,.omega.-diorganopolysiloxane diol, organohydrogenpolysiloxane, quaternary ammonium salt catalyst and a heavy metal carboxylate.\nIn Nippon Tokkyo Koukoku Koho 45-11839 (Japanese Patent Publication) Murphy invented an RTV foam-forming organopolysiloxane composition comprising .alpha.,.omega.-diorganopolysiloxane diol, organohydrogenpolysiloxane and aminohydroxyorganopolysiloxane compounds as the dehydrogenation catalyst.\nIn Nippon Tokkyo Kokai Koho 51-46352 Endo invented RTV foam-forming organopolysiloxane compositions comprising .alpha.,.omega.-diorganopolysiloxane diol, organohydrogenpolysiloxane and a platinum catalyst.\nHowever, all of these technologies gave RTV foam-forming compositions having low foam-forming capability and inferior foam-formability. Accordingly, the previous compositions could only be applied to limited commercial production. The present invention is based on the discovery of a new class of RTV foam-forming compositions which overcome the defects of previous technologies and provide improved compositions having high foam-forming ratios."}
{"text": "1. Field of the Invention\nThe present invention relates to a process and devices for detecting the instant of injection and for determining the duration of injection in hemodynamic monitoring by means of thermodilution.\n2. Description of the Prior Art\nThe measurement of hemodynamic parameters, for example the cardiac output, is largely performed at present either by means of pulmonary arterial or transcardiopulmonary thermodilution (Pfeiffer U. J., Knoll R. (1993): Process for Determining a Patient's Circulatory Fill Status. U.S. Pat. No. 5,526,817) or else by means of thermo-dye-dilution (Pfeiffer, U. J., Backus G., Blumel G., Eckart J., Muller P., Winkler P., Zeravik J., Zimmermann G. J. (1990): A Fiberoptics-Based System for Integrated Monitoring of Cardiac Output, Intrathoracic Blood Volume. Extravascular Lung Water, O.sub.2 Saturation, and a-v Differences. Practical Applications of Fiberoptics in Critical Care Monitoring, Springer Verlag, 114-125). In these processes, a defined volume of an indicator substance which is as cold as possible, for example glucose or saline solution, is injected. The instant of injection into the body is registered by means of an extracorporeal temperature sensor which is integrated directly in the injection lumen.\nAt the same time, the thermodilution measurement is started by means of a thermosensor, which in the case of pulmonary arterial measurement is located in the distal lumen of the pulmonary artery catheter in the Arteria pulmonalis or, in the case of transcardiopulmonary measurement, in the tip of a catheter lying in the Arteria femoralis or in the Aorta abdominalis. By plotting the thermodilution curve, the cardiac output can be calculated, for example by means of the Stewart-Hamilton method.\nThe special aspect of the transcardiopulmonary method is the additional determination of a number of cardiovascular parameters, in particular for assessing the output status, for example by the intrathoracic blood volume. For the calculation of these parameters, knowledge of the characteristic times of the indicators, in particular the mean transit time and exponential fall time, is required. To be able to calculate these exactly, the instant of injection, the mean passage time of the injectate and the duration of injection must in turn be accurately measured, which is accomplished by means of the curve plotted using the extracorporeal temperature sensor (cf. FIG. 1, which reproduces the injection curve profile with a known injectate temperature sensor system; in contrast to this, FIG. 2 shows the injection curve profile with a sensor system according to the invention set out below).\nAs a function of the temperature difference between ambient air and injectate, the value T.sub.inj, required for correct measurements, is calculated using additionally determined correction factors.\nA major disadvantage of the existing technique is that injectate of a temperature deviating from room temperature was required for optimum measurements in order to determine exactly the instant of injection and the duration of injection, since the volume in the customary extracorporeal injectate temperature sensor housing is essentially at room temperature. To be able to detect the instant at which injection starts and to be able to calculate the duration of injection from the temperature profile, a clear temperature difference between the fluid at the sensor before injection and the injectate is required.\nFor this reason, it must be ensured that the injection solution is available in a well cooled state at any time. This means additional work also for the nursing staff in intensive care units and in operating rooms. In addition, measurements often do not proceed absolutely smoothly, with the result that injectate taken out of cooling too early may already have warmed up again by the time it is used; the same problem arises if a number of measurements are carried out at short intervals one after the other.\nThe use of cooling sets, which can be installed at the patient's bed, does offer the advantage of an injectate cooled for a certain time directly at the patient, but again brings about considerable disadvantages due to increased work, for example to obtain fresh ice for the cooling box, and due to the costs additionally incurred."}
{"text": "A lithography technique is a very important process among semiconductor manufacturing processes by which scaling-down of a semiconductor device is achieved, because the lithography technique is a process that generates a pattern of the device. Recently, according to high integrity of LSI, circuit line width required for a semiconductor device has been reduced year after year. In order to form a desired circuit pattern on the semiconductor device, a highly-accurate master image pattern (sometimes, referred to as a reticle or a mask) is needed. Herein, since an electron beam writing technique intrinsically has excellent resolution, the technique is used to produce a highly-accurate master image pattern.\nIn the above-described electron beam writing, uniformity of line width in more accurate sample plane, for example, a mask plane is required. Herein, in the electron beam writing, electrons are charged in a deflector, and thus, the electron beam is drifted, so that there occurs a phenomenon that position accuracy of the writing is degraded.\nIn order to improve the position accuracy of the writing, it is preferable that the drift of the electron beam is suppressed.\nJP-A H06-120126 discloses a technique of manufacturing an aperture plate by using tungsten having high electron beam blocking ability so as to improve processing accuracy of an opening portion of the aperture plate."}
{"text": "The technical field of this invention concerns medical devices useful for the repair of injured nerves and methods for preparing and using such devices for nerve repairs.\nThe problem of repairing severed nerves is a long-standing one that has plagued surgeons for over a hundred years. Despite advances in microsurgical techniques, a patient's recovery from a serious wound is often limited by a degree of nerve damage which cannot be repaired. The replanting of amputated fingers and limbs is especially limited by poor nerve regeneration.\nWhen a nerve is severed, the functions supplied by that nerve, both motor and sensory, are lost. The nerve cells' appendages (axons) in the distal (the furthest away from the spinal cord) portions of the severed nerve degenerate and die leaving only the sheaths in which they were contained. The axons in the proximal stump that are still connected to the spinal cord or dorsal root ganglion, also suffer some degeneration. The degeneration generally does not proceed to the death of the entire nerve cell bodies. If the injury occurs far enough from the nerve cell bodies, regeneration will occur. Axonal sprouts will appear from the tip of the regenerating axon. These sprouts grow distally and attempt to reenter the intact neurilemnal sheaths of the distal portion of the severed nerve. If entry is successfully made, axonal growth will continue down these sheaths and function will eventually be restored.\nIn the conventional approach to nerve repair, an attempt is made to align the cut ends of the fascicles (nerve bundles within the nerve trunk). A similar approach is taken with smaller nerves. In either case, the chief hazard to the successful repair is the trauma produced by the manipulation of the nerve ends and the subsequent suturing to maintain alignment. The trauma appears to stimulate the growth and/or migration of fibroblasts and other scar-forming connective tissue cells. The scar tissue prevents the regenerating axons in the proximal stump from reaching the distal stump to reestablish a continuous pathway. The result is a permanent loss of sensory or motor function.\nVarious attempts have been made over the years to find a replacement for direct (i.e., nerve stump-to-nerve-stump suturing). Much of the research in this field has focused on the use of \"channels\" or tubular prostheses which permit the cut ends of the nerve to be gently drawn into proximity and secured in place without undue trauma. It is also generally believed that such channels can also prevent, or at least retard, the infiltration of scar-forming connective tissue.\nThe use of silastic cuffs for peripheral nerve repair was reported by Ducker et al. in Vol. 28, Journal of Neurosurgery, pp. 582-587 (1968). Silicone rubber sheathing for nerve repair was reported by Midgley et al. in Vol. 19, Surgical Forum, pp. 519-528 (1968) and by Lundborg, et al. in Vol. 41, Journal of Neuropathology in Experimental Neurology, pp. 412-422 (1982). The use of bioresorbable polyglactin mesh tubing was reported by Molander et al. in Vol. 5, Muscle & Nerve, pp. 54-58 (1982). The use of semipermeable acrylic copolymer tubes in nerve regeneration was disclosed by Uzman et al. in Vol. 9, Journal of Neuroscience Research, pp. 325-338 (1983). Bioresorbable nerve guidance channels of polyesters and other polymers have been reported by Nyilas et al. in Vol. 29, Transactions Am. Soc. Artif. Internal Organs, pp. 307-313 (1983) and in U.S. Pat. No. 4,534,349 issued to Barrows in 1985.\nDespite the indentification of various materials which can serve as nerve guidance channels, the results of research to date have revealed significant shortcomings in such prostheses. Some of the materials identified above have lead to inflammatory reactions in the test animals and have failed to exclude scar tissue formation within the channels. Moreover, the total number of axons, the number of myelinated axons, the thickness of the epineurium, and the fascicular organization of nerves regenerated within guidance channels are all typically less than satisfactory and compare poorly with the original nerve structure of the test animals. Moreover, the loss of sensory or motor function is still the most common outcome of such laboratory experiments.\nThere exists a need for a better materials and methods for formation of nerve guidance channels. Materials and methods for nerve repair that would minimize surgical trauma, prevent interference with nerve growth by scar tissue, and improve the chances for successful recovery of sensory or motor function, would satisfy a long-felt need in this field."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of compressed air and gas systems, and more particularly to air/gas-driven tools, such as paint guns and other equipment. The present invention further relates to dryness indicators designed to visually indicate the moisture content of compressed gas or air delivered to a point of use.\n2. Description of Prior Art\nTypical compressed air or gas produced by a compressor apparatus is saturated with 50% to 100% relative humidity. Removal of this moisture vapor requires that an air/gas drying system be used, such as a refrigerated dryer or an adsorbent type of dryer. Such apparatus are generally very effective, the latter being typically capable of drying compressed air or gas to below-zero dew point levels. Notwithstanding such drying measures, there is unfortunately no guarantee that the compressed air or gas will have the desired dryness by the time it arrives through an air/gas feed system to a downstream point of use. Compressed air lines, various fitting and regulation devices, or improper operation of the dryer system all represent sources of residual moisture vapor entrainment in the air/gas feed system. This means that moisture-treated compressed air or gas may be carrying unwanted moisture vapor when it goes into use as an application.\nOne area where this problem tends to occur is in paint booth operations where compressed air or gas is used as a propellant to atomize and expel paint from a paint gun. Even though extraordinary measures are often implemented to eliminate moisture vapor at the compressed air/gas source, moisture can still be delivered to the paint gun. In some cases, this may be due to the drying system losing effectiveness due to a malfunction or other problem. However, even if the drying system is operating at full operational efficiency, the lengthy hoses connecting the air source to the paint gun can introduce unwanted moisture vapor into the system. In particular, these hoses can be disconnected and re-connected any number of times throughout the course of a painting application. Each time a disconnection occurs, moisture-laden ambient air is allowed to enter the air/gas line, and will feed through the paint gun until such time as it evacuated from the line and replaced by dry air/gas coming from the air/gas source. Any time there is excess moisture vapor in a paint gun, unwanted fouling can occur that results in a bad and unacceptable paint job. In most cases, the unsuspecting painter will assume the air/gas quality is satisfactory, particularly when there is sophisticated drying equipment operating at the air/gas source.\nVarious dryness indicators have been proposed for use in compressed air, gas and refrigerant applications. These typically involve the use of a moisture-adsorbing silica gel desiccant that is impregnated with a chemical moisture indicator, such as a cobalt salt. This particular chemical indicator is normally a deep blue color when it is dry, but gradually turns a light pink color in proportion to the amount of moisture that is present as the salt hydrates. In a dryness indicator, the color-indicating desiccant is placed in contact with a compressed air/gas stream within a transparent or translucent container, so that the desiccant can be viewed during operations.\nPrior art dryness indicators tend to have design features that prevent them from being optimally suited for point-of-use operation in conjunction with a hand-held air/gas-driven tool, such as a paint gun, where moisture monitoring is most needed. In all of the reference materials reviewed, the prior art dryness indicators form part of a filter/dryer that requires a relatively large quantity of desiccant to effectively remove moisture for a reasonable length of time. This quantity of desiccant is more than that which is required to indicate dryness. The filter/dryers in which prior art dryness indicators are incorporated also tend to include additional elements to condense and remove moisture droplets from the air/gas stream, and to trap oil, line debris and other contaminants. As a result of the foregoing design features, most prior art dryness indicators are large or bulky, and not suitable for point-of-use operation.\nIt is to solving the foregoing problems that the present invention is directed. What is particularly needed is an improved compressed air/gas dryness indicator that is optimized for point-of-use operation with an air/gas-driven tool, such as a paint gun. Ideally, the dryness indicator needs to provide a visual indication identifying the exact state of dryness of the compressed air/gas line, yet must be unobtrusive and afford full freedom of movement at the point-of-use without any impediment of bulky filters, desiccant containers, cumbersome vessels, etc. The dryness indicator additionally needs to be easy to install and use, should be simple and inexpensive, and should require little or no maintenance."}
{"text": "The present invention relates generally to the imaging of objects in highly scattering turbid media and more particularly to a novel technique for imaging objects in highly scattering turbid media.\nAs can readily be appreciated, there are many situations in which the detection of an object present in a highly scattering turbid medium is highly desirable. For instance, the detection of a tumor embedded within a tissue is one such example. Although X-ray techniques do provide some measure of success in detecting objects in turbid media, they are not well-suited for detecting very small objects, e.g., tumors less than 1 mm in size, or for detecting objects in thick media. In addition, X-ray radiation can present safety hazards to a person exposed thereto.\nAn alternative technique used to detect objects in turbid media is transillumination. In transillumination, visible or near infrared (NIR) light is incident on one side of a medium and the light emergent from the opposite side of the medium is used to form an image. Objects embedded in the medium typically absorb the incident light and appear in the image as shadows. Unfortunately, the usefulness of transillumination as a detection technique is severely limited in those instances in which the medium is thick or the object is very small. This is because light scattering within the medium contributes to noise and reduces the intensity of the unscattered light used to form the image shadow.\nTo improve the detectability of small objects located in a turbid medium using transillumination, many investigators have attempted to selectively use only certain components of the transilluminating light signal. This may be done by exploiting the properties of photon migration through a scattering medium. Photons migrating through a turbid medium have traditionally been categorized into three major signal components: (1) the ballistic (coherent) photons which arrive first by traveling over the shortest, most direct path; (2) the snake (quasi-coherent) photons which arrive within the first .delta.t after the ballistic photons and which deviate, only to a very slight extent, off a straight-line propagation path; and (3) the diffusive (incoherent) photons which experience comparatively more scattering than do ballistic and snake photons and, therefore, deviate more considerably from the straight-line propagation path followed by ballistic and snake photons.\nBecause it has been believed that ballistic and snake photons contain the least distorted image information and that diffusive photons lose most of the image information, efforts to make transillumination work most effectively with turbid media have focused on techniques which permit the selective detection of ballistic and snake photons while rejecting diffusive photons. This process of selection and rejection has been implemented in various time-gating, space-gating and time/space-gating techniques. Patents, patent applications and publications which disclose certain of these techniques include U.S. Pat. No. 5,140,463, inventors Yoo et al., which issued Aug. 18, 1992; U.S. Pat. No. 5,143,372, inventors Alfano et al., which issued Aug. 25, 1992; U.S. Pat. No. 5,227,912, inventors Ho et al., which issued Jul. 13, 1993; U.S. Pat. No. 5,371,368, inventors Alfano et al., issued Dec. 6, 1994; Alfano et al., \"Photons for prompt tumor detection,\" Physics World, pp. 37-40 (January 1992); Wang et al., \"Ballistic 2-D Imaging Through Scattering Walls Using an Ultrafast Optical Kerr Gate,\" Science, Vol. 253, pp. 769-771 (Aug. 16, 1991); Wang et al., \"Kerr-Fourier imaging of hidden objects in thick turbid media,\" Optics Letters, Vol. 18, No. 3, pp. 241-243 (Feb. 1, 1993); Yoo et al., \"Time-resolved coherent and incoherent components of forward light scattering in random media,\" Optics Letters, Vol. 15, No. 6, pp. 320-322 (Mar. 15, 1990); Das et al., \"Ultrafast time-gated imaging in thick tissues: a step toward optical mammography,\" Optics Letters, 18(13):1092-4 (1993); Chen et al., \"Two-dimensional imaging through diffusing media using 150-fs gated electronic holography techniques,\" Optics Letters, Vol. 16, No. 7, pp. 487-489 (Apr. 1, 1991); Duncan et al., \"Time-gated imaging through scattering media using stimulated Raman amplification,\" Optics Letters, Vol. 16, No. 23, pp. 1868-1870 (Dec. 1, 1991), all of which are incorporated herein by reference.\nOf the above-listed art, Wang et al., \"Kerr-Fourier imaging of hidden objects in thick turbid media,\" Optics Letters, Vol. 18, No. 3, pp. 241-243 (Feb. 1, 1993) is illustrative of transillumination techniques which selectively use the ballistic and/or snake components of light. In this article, there is disclosed a time/space-gating system for use in imaging opaque test bars hidden inside a 5.5 cm-thick 2.5% Intralipid solution. The disclosed system includes three main parts: a laser source, an optical Kerr gate and a detector. The laser source is a picosecond mode-locked laser system, which emits a 1054 nm, 8 ps laser pulse train as the illumination source. The second harmonic of the pulse train, which is generated by transmission through a potassium dihydrate phosphate (KDP) crystal, is used as the gating source. The illumination source is sent through a variable time-delay and is then used to transilluminate, from one side, the turbid medium containing the opaque object. The signal from the turbid medium located at the front focal plane of a lens is collected and transformed to a Kerr cell located at its back focal plane (i.e., the Fourier-transform spectral plane of a 4F system). That portion of the Kerr cell located at the focal point of the 4F system is gated at the appropriate time using the gating source so that only the ballistic and snake components are permitted to pass therethrough. The spatial-filtered and temporal-segmented signal is then imaged by a second lens onto a CCD camera.\nAlthough techniques of the type described above, which selectively use ballistic and snake photons to image objects in turbid media, have enjoyed a modicum of success, such techniques have been limited by the fact that detected light signals derived from ballistic and snake photons are typically rather weak, due to the proportionately small number of transilluminated ballistic and snake photons. This problem is further exacerbated in those instances in which the turbid medium is thick and the likelihood of substantial scattering increases.\nAccordingly, because diffusive photons constitute the greatest component of the transilluminated light signal, it would be highly desirable to make use of the diffusive component of the light signal in forming an image via transillumination. This objective is made difficult, however, by the fact that diffusive photons tend to traverse a medium along ill-defined paths. One approach to this problem has been to invert the experimental scattering data obtained from various points in the medium using some inverse algorithm and reconstruction approach. This approach is often called diffusion tomography since diffusion or scattering is the dominant factor in the problem. In diffusion tomography, one produces an internal map of the scattering medium using the scattered signals and a mathematical inversion algorithm. The inversion is based upon the physical and mathematical principles governing photon propagation in turbid media. Both time-resolved data and frequency domain data can be used for reconstruction. Examples of diffusion tomography techniques include Arridge, \"The Forward and Inverse Problems in Time Resolved Infra-Red Imaging,\" Medical Optical Tomography: Functional Imaging and Monitoring SPIE Institutes, Vol. IS11, G. Muller ed., 31-64 (1993); Singer et al., \"Image Reconstruction of the Interior of Bodies That Diffuse Radiation,\" Science, 248:990-3 (1993); Barbour et al., \"A Perturbation Approach for Optical Diffusion Tomography Using Continuous-Wave and Time-Resolved Data,\" Medical Optical Tomography: Functional Imaging and Monitoring SPIE Institutes, Vol. IS11, G. Muller ed., 87-120 (1993); M. Patterson et al., SPIE, 1767, 372 (1992); J. Schotland et al., App. Opt., 32, 448 (1993), all of which are incorporated herein by reference.\nThe foregoing diffusion tomography techniques do not lead to a resolution that is better than about 5-10 mm. Moreover, these techniques are time-consuming and do not readily lend themselves to real-time use."}
{"text": "This invention is concerned with apparatus for a ball game."}
{"text": "This invention relates to an apparatus and method for the production of hydrocarbons from earth formations, and more particularly, to those hydrocarbon-bearing deposits where the oil viscosity and saturation are so high that sufficient steam injectivity cannot be obtained by current steam injection methods. Most particularly this invention relates to an apparatus and method for the production of hydrocarbons from tar sand deposits containing layers of high electrical conductivity and having vertical hydraulic connectivity between the various geologic sequences.\nReservoirs in many parts of the world are abundant in heavy oil and tar sands. For example, those in Alberta, Canada; Utah and California in the United States; the Orinoco Belt of Venezuela; and the USSR. Such tar sand deposits contain an energy potential estimated to be quite great, with the total world reserve of tar sand deposits estimated to be 2,100 billion barrels of oil, of which about 980 billion are located in Alberta, Canada, and of which 18 billion barrels of oil are present in shallow deposits in the United States.\nConventional recovery of hydrocarbons from heavy oil deposits is generally accomplished by steam injection to swell and lower the viscosity of the crude to the point where it can be pushed toward the production wells. In those reservoirs where steam injectivity is high enough, this is a very efficient means of heating and producing the formation. Unfortunately, a large number of reservoirs contain tar of sufficiently high viscosity and saturation that initial steam injectivity is severely limited, so that even with a number of \"huff-and-puff\" pressure cycles, very little steam can be injected into the deposit without exceeding the formation fracturing pressure. Most of these tar sand deposits have previously not been capable of economic production.\nIn steam flooding deposits with low injectivity the major hurdle to production is establishing and maintaining a flow channel between injection and production wells. Several proposals have been made to provide horizontal wells or conduits within a tar sand deposit to deliver hot fluids such as steam into the deposit, thereby heating and reducing the viscosity of the bitumen in tar sands adjacent to the horizontal well or conduit. U.S. Pat. No. 3,986,557 discloses use of such a conduit with a perforated section to allow entry of steam into, and drainage of mobilized tar out of, the tar sand deposit. U.S. Pat. Nos. 3,994,340 and 4,037,658 disclose use of such conduits or wells simply to heat an adjacent portion of deposit, thereby allowing injection of steam into the mobilized portions of the tar sand deposit.\nSeveral prior art proposals designed to overcome the steam injectivity problem have been made for various means of electrical or electromagnetic heating of tar sands. One category of such proposals has involved the placement of electrodes in conventional injection and production wells between which an electric current is passed to heat the formation and mobilize the tar. This concept is disclosed in U.S. Pat. Nos. 3,848,671 and 3,958,636. A similar concept has been presented by Towson at the Second International Conference on Heavy Crude and Tar Sand (UNITAR/UNDP Information Center, Caracas, Venezuela, Sept. 1982). A novel variation, employing aquifers above and below a viscous hydrocarbon-bearing formation, is disclosed in U.S. Pat. No. 4,612,988. In U.S. Pat. No. Re. 30738, Bridges and Taflove disclose a system and method for in-situ heat processing of hydrocarbonaceous earth formations utilizing a plurality of elongated electrodes inserted in the formation and bounding a particular volume of a formation. A radio frequency electrical field is used to dielectrically heat the deposit. The electrode array is designed to generate uniform controlled heating throughout the bounded volume.\nIn U.S. Pat. No. 4,545,435, Bridges and Taflove again disclose a waveguide structure bounding a particular volume of earth formation. The waveguide is formed of rows of elongated electrodes in a \"dense array\" defined such that the spacing between rows is greater than the distance between electrodes in a row. In order to prevent vaporization of water at the electrodes, at least two adjacent rows of electrodes are kept at the same potential. The block of the formation between these equipotential rows is not heated electrically and acts as a heat sink for the electrodes. Electrical power is supplied at a relatively low frequency (60 Hz or below) and heating is by electric conduction rather than dielectric displacement currents. The temperature at the electrodes is controlled below the vaporization point of water to maintain an electrically conducting path between the electrodes and the formation. Again, the \"dense array\" of electrodes is designed to generate relatively uniform heating throughout the bounded volume.\nHiebert et al (\"Numerical Simulation Results for the Electrical Heating of Athabasca Oil Sand Formations,\" Reservoir Engineering Journal, Society of Petroleum Engineers, Jan. 1986) focus on the effect of electrode placement on the electric heating process. They depict the oil or tar sand as a highly resistive material interspersed with conductive water sands and shale layers. Hiebert et al propose to use the adjacent cap and base rocks (relatively thick, conductive water sands and shales) as an extended electrode sandwich to uniformly heat the oil sand formation from above and below.\nThese examples show that previous proposals have concentrated on achieving substantially uniform heating in a block of a formation so as to avoid overheating selected intervals. The common conception is that it is wasteful and uneconomic to generate nonuniform electric heating in the deposit. The electrode array utilized by prior inventors therefore bounds a particular volume of earth formation in order to achieve this uniform heating. However, the process of uniformly heating a block of tar sands by electrical means is extremely uneconomic. Since conversion of fossil fuel energy to electrical power is only about 38 percent efficient, a significant energy loss occurs in heating an entire tar sand deposit with electrical energy.\nGeologic conditions can also hinder heating and production. For example, many formations have little or no vertical hydraulic connectivity within the formation. This means that once the selected layer is preheated, vertical movement of the steam will be somewhat limited, thus limiting vertical transfer of heat to that which can be carried by thermal conduction. However, in other instances, the geologic conditions can actually help production, provided that the recovery method is designed to take advantage of the geologic conditions. In formations in which there is vertical hydraulic connectivity, once steam is injected into a layer, the heated oil progressively drains downwards within the deposit, allowing the steam to rise within the deposit. The steam flowing into the tar sand deposit effectively displaces oil toward the production wells, and provides heat to the formation.\nU.S. Pat. No. 4,926,941 (Glandt et al) discloses electrical preheating of a thin layer by contacting the thin layer with a multiplicity of vertical electrodes spaced along the layer.\nIt is therefore an object of this invention to provide an efficient and economic method of in-situ heat processing of tar sand and other heavy oil deposits having vertical hydraulic connectivity, wherein electrical current is used to heat thin layers within such deposits, utilizing a minimum of electrical energy to prepare the tar sands for production by steam injection; and then to efficiently utilize steam injection to mobilize and recover a substantial portion of the heavy oil and tar contained in the deposit."}
{"text": "A typical computer-aided design (CAD) system employed in engineering contexts uses a geometry-oriented approach to define and represent engineering information. As a result, a designer is required to perform a number of low-level geometric operations in order to produce a final digital model of a desired product. In other words, the designer has to focus on the details of geometry creation rather than the required functionality of the product.\nIn addition, conventional CAD systems use a rigid history-based modeling approach that creates dependency between the operations performed by the designer when creating a digital model. That is, an earlier operation may influence subsequent operations or, alternatively, a subsequent operation may obliterate the functionality of a prior operation. This order dependent nature of feature operations makes it difficult for the designer to apply any modification to the digital model or correct mistakes occurred early in the design process.\nFurther, conventional CAD systems do not allow members of a product development team to work together simultaneously, regardless of their location or role, in order to create and finalize product definitions. Instead, current design intent collaborative practices are based on serial collaboration, in which copies of digital product files are passed back and forth among development team members who must wait to get these copies back before they can make or even suggest design refinements."}
{"text": "1. Field of the Invention\nThe present invention relates generally to securing information in computing systems, and more specifically to limiting access to that information based on the context in which at least a portion of the transactional information was generated, such as from a sale.\n2. Background\nSecuring information has become a priority for organizations to ensure that business processes and information relating thereto remain confidential. As an organization's business processes becomes more complex, the means for securing information has to be flexible to adapt to organizational changes while preserving an appropriate balance between confidentiality (i.e., limiting access to information) and openness (i.e., freedom to access information), both of which are necessary for the success of the organization. Examples of organizational changes requiring such flexibility include employee/group transfers, company reorganizations, compensation plan adjustments and the hiring and/or terminating of personnel.\nTo manage compensation schemes through these types of organizational changes, as well as providing incentive-based compensation for employees in general, organizations have structured compensation plans in accordance with Enterprise Incentive Management (EIM) principles. These principles tailor compensation plans so as to improve optimal performance and to align the organization's strategy with the desired behaviors of it employees. EIM refers generally to managing variable pay plans throughout an organization (i.e., corporation or enterprise) and includes plans for salespeople, suppliers, distribution channel partners, brokers, customers, employees, executives, and partners.\nBut conventional approaches to securing information generated in the framework of an organization typically lack the flexibility to adapt to changes in corporate processes or structure, such as a change in traditional compensation schemes or personnel. For example, consider a personnel change from one part of an organization to another part as shown in FIG. 1.\nFIG. 1 depicts a traditional organizational chart illustrating an employee transferring from one position in organizational structure 100 to another position in new organizational structure 110. Organizational structures 100 and 110 each represent a hierarchical structure depicting supervisor-subordinate relationships where permissions to access secured information decreases from the top position occupied by “A” to the bottom positions occupied by “D,” “E,” and “F.” A square box in FIG. 1, such as the one labeled A, represents a position or role occupied by an employee (or a group/organizational element) and is interrelated with other square boxes, where the interrelations are depicted as lines connecting at least two square boxes. Hence, an employee or organizational element occupying box A is in a supervisory role to employees or organizational elements in boxes “B” and “Cynthia,” which are both in subordinate roles to that of box A.\nIn organizational structure 100, Cynthia is shown to occupy a supervisory role in relation to boxes E and F, which can be employees, groups of employees or other organizational elements. In this role, Cynthia has a “span of access” 102 and is granted permission to access information relating at least to her subordinates occupying boxes E and F, which may include transactional information forming Cynthia's compensation.\nFurther, consider that an employee associated with box E is a sales person operating according to a compensation plan that specifies each of the following allocations to their supervisor Cynthia's compensation: 2% of sales revenue within a particular geographic region; 1% of sales from a particular product line; 2% of sales to a particular customer; and 0.5% of sales by other members of her sales team. Since Cynthia's compensation is based upon such a compensation scheme, she and others in her position are traditionally authorized to access transactional information in her span of access 102 to verify that the sales person's sales revenues are accurately recorded. In particular, Cynthia can access transactional information for boxes E and F but not boxes B and D, which are outside of her span of access 102.\nNext, consider that Cynthia assumes a new role in new organizational structure 110 in the position formerly demarcated as box B of organization structure 100. In this role, traditional security mechanisms allow Cynthia to now have access to transactional information within span of access 112, which authorizes her to examine the activities of the employee(s) relating to box D to review the transactional information that affects her compensation in this new role. But once Cynthia assumes this new role in organizational structure 110, she traditionally is precluded from having access to transactional information for span of access 102. This is generally due to traditional approaches to securing information where a person's set of permissions are dependent on the person's current position in an organization. Without having span of access 102, Cynthia is unable to determine whether her previous efforts and those of her previous subordinates are adequately and accurately recognized so that she can justly be compensated for any activity occurring before she assumed her new role in organization structure 110. Thus, there is a need to provide a flexible method of securing information such that the aforementioned drawbacks of conventional EIM schemes are overcome."}
{"text": "1. Field of the Invention\nThe present invention relates to a vehicle drive assist system and more particularly to a vehicle drive assist system capable of detecting an oncoming vehicle which is difficult to be recognized at intersections based on frontal information obtained from a stereoscopic camera and the like.\n2. Discussion About Prior Arts\nWhen a vehicle make a right turn at intersections in the “keep to the left” traffic system, in case where there are oncoming vehicles waiting for turning right on oncoming lanes, it is difficult for a driver to confirm oncoming vehicles traveling straight. Under such situations, the driver must put miscellaneous information such as oncoming vehicles, pedestrians walking across a road ahead and the like in order and therefore he or she is forced to bear lots of burdens.\nIn order to reduce such burdens of the driver, Japanese Patent Application Laid-open No. Toku-Kai-Hei 9-282592 discloses a technique in which a collision of a vehicle turning right with an oncoming vehicle traveling straight in intersections is prevented by detecting the oncoming vehicle with an obstacle sensor installed in intersections and warning a driver.\nHowever, the obstacle sensor and the warning system must be installed in every intersection and a huge amount of money is needed.\nFurther, Japanese Patent Application Laid-open No. Toku-Kai 2001-101592 discloses a technique wherein a vehicle itself has an ability to detect oncoming vehicles at an early stage using a fish-eye lens installed in a bumper of the vehicle.\nHowever, this technique has such problems that stains, raindrops and the like sticking to the lens hinder accurate imaging and also images taken through the fish-eye lens require complicated correction processing."}
{"text": "Complex global supply chains for mobile phones, computers, printers, automobiles, aircraft, defense systems, medical equipment and other important products are subject to risks and vulnerabilities that enable infiltration of counterfeit goods into legitimate trade channels. Counterfeiting of industrial and consumer products, and components incorporated into these products, can compromise the integrity of final products, generate losses to legitimate businesses and expose consumers to fake, faulty or harmful products.\nExisting, widely deployed anti-counterfeiting technologies are subject to various limitations. For example one-dimensional (1D) and two-dimensional (2D) bar codes, holograms, and other optically readable labels and tags themselves be counterfeited (i.e., copied). Radio frequency identification (RFID) tags have raised privacy concerns and also incur a relatively high unit cost, which generally makes them uneconomical for use with high volume, low unit cost products such as active and passive electronic components."}
{"text": "1. Field of the Invention\nThe invention relates to ski bindings used to hold boots onto skis.\n2. Description of Prior Art\nGerman Pat. No. 222,828 discloses a ski binding for maintaining one end of a boot onto a ski."}
{"text": "The present invention relates to methods of forming precision metal parts and, more specifically, to thixotropic forming of precision multi-alloy parts.\nAs performance criteria for turbine engines becomes more stringent, there is a need for an improved turbine rotor that exhibits maximum resistance to both fatigue and creep.\nDie casting is a well-known process for producing complex components with excellent surface quality and good dimensional accuracy. However, the structural integrity of die castings is often compromised by air trapped in the casting upon injection of the liquid metal into the die casting cavity. The resultant porosity also compromises heat treatment of the casting that is often necessary to refine the grain structure and increase the strength of the casting.\nForging is also a well known process for producing relatively strong components having a desirable grain structure. However, forged products generally exhibit relatively low resistance to creep due to their fine grain structures.\nThixotropic, or semisolid, metal forming is a viable alternative to traditional casting and forging methods. This process lies somewhere between a casting and a forging process in that the metal to be formed is brought to a xe2x80x9cthixotropicxe2x80x9d state; that is, 30 or 40 percent of the mass exists in a liquid phase and the balance in a solid phase. The solid portion comprises small spherically-shaped nodules suspended within the liquid phase. Semisolid metals heated to a thixotropic state exhibit unique Theological properties due to their non-dendritic, or spherical, microstructure. The rheological properties of the semisolid metal range from high viscosities, like table butter, for alloys at rest, to low viscosities, such as machine oil, as the shearing rate of the semisolid slug is increased. By heating the metals to a semisolid range and then agitating the semisolid alloy, the dendritic microstructure normally found is eliminated and replaced by the spherical microstructure. Upon solidification, the alloys then exhibit a fine equiaxed microstructure.\nNormally, a highly viscous thixotropic slug will retain its outer shape provided there are no external forces, other than gravity, applied to it. However, its butter-like consistency is easily deformed to a low viscosity, particularly by a shearing action such as high velocity impact, making it extremely suitable when driving the alloy into the mold during the manufacturing process. Because semisolid-formed alloys exhibit an intermediate-sized grain structure, larger than forged grains and smaller than cast grains, it is expected that semisolid forged or cast alloys will have improved creep rupture resistance over traditionally forged alloys and improved strength properties over traditionally cast alloys.\nThe thixotropic process has been extensively studied by others in relation to lighter metals or low melting point metals such as aluminum, magnesium, zinc, and copper alloys. In general, the low melting point metals have a melting point within the range of 750 to 1250 degrees Fahrenheit. On the other hand, lower melting point copper alloys (other than Cuxe2x80x94Ni alloys for example), have a melting point from 1200-1900 degrees Fahrenheit but are still within the scope of the present invention.\nVery little research is available with regard to high temperature alloys commonly used in turbine rotors, including ferrous or nickel-based alloys. In general, high temperature alloys have a melting point range from 2000xc2x0 F.-2700xc2x0 F. One significant difference between semisolid production for lighter alloys and that for high temperature alloys involves the adaptation of the process to the problematic and high heating temperatures of 2500xc2x0 F. to 2700xc2x0 F. as opposed to alloys in the 750xc2x0 F.-1250xc2x0 F. melting point range. Designing a semisolid process compatible with such high heat has proven challenging. Generally, chrome-nickel alloys of, for example, 18% Cr and 82% Ni are used in turbine rotor forgings. This alloy has a solidus of 2550xc2x0 F., and a liquidus of 2640xc2x0 F. where the alloy is completely molten. The semisolid/thixotropic phase exists between the solidus and liquidus temperatures at temperatures ranging between 2550xc2x0 F. and 2640xc2x0 F. The alloy is commonly forged at temperatures below 2550xc2x0 F., in the solid phase, and cast at molten temperatures above 2640xc2x0 F., in the liquid phase.\nYet another problem is that current net-shape forging and die-casting equipment design includes permanent molds that often do not readily separate from the part interface when removing the turbine rotors and their intricate blades from the mold. This results in fractured or weakened blades and a corresponding number of rejected parts that do not meet design specifications. A need exists for semisolid manufacturing methods that facilitate ease of removal of the finished part, thereby improving the production volume and reducing the rejection rate of the finished parts.\nFinally, precision metal assemblies are specifically designed to withstand various forces under uniquely stressful conditions. In certain applications, however, one part of a complete assembly may be exposed to stress and temperature loads significantly different from that of other parts integral to the same assembly. For example, the bore of a rotor may require good elongation, high strength, and good low cycle fatigue properties but may not require high temperature properties. In contrast, certain blade or rim portions of the rotor might require very high creep resistance and stress rupture strength at elevated temperatures. Formulating a single alloy capable of withstanding the variable stresses subjected to different locations within a precision metal assembly has also proven challenging.\nEuropean Pat. No. 0 574 141 A1 entitled, xe2x80x9cThixoformable Layered Materials and Articles Made From Themxe2x80x9d, discloses a method of sequentially applying layers of substantially metallic material in a thixoforming process. The reference discloses a rotatable cylindrical collector that collects molten metal cooled by inert gases prior to concentric deposition on the collector. A thixotropic layer comprising 30-70% liquid is thereby formed as the atomized metal is sprayed onto the collector. Additional layers are then added in the same way that may or may not comprise the same alloy as found in the first layer. At least two of the layers have different properties. The layers may also comprise reinforcing materials such as ceramic, metallic, and intermetallic materials in spherical, fibrous, or any other shape. The reinforcing materials may be added by simply spraying them into the atomized melt spray during that stage.\nEP 0 574 141 A1 discloses that a layered composition thixoformed by this method exhibits enhanced toughness and damage resistance due to the layered 3-dimensional structure. However, the method is labor intensive, for each layer must at least be melted, sprayed, cooled by inert gases, and then collected on the surface of the cylinder. The preforms formed by this method must then be cooled and cut to accommodate a thixoforming forging process, wherein a multiproperty component is manufactured.\nU.S. Pat. Nos. 5,832,982, 5,878,804, and 6,003,585 to Williams et al. describe metal forming processes that form a semisolid slug or thixotropic solution containing 30-40% liquid and 60-70% solids. The slug is inductively heated to destroy the dendritic microstructure and when hardened results in a preferred fine equiaxed microstructure. A disadvantage of the processes, however, is that the slug must be heated to a point wherein the liquid percentage exceeds 40%, thereby ensuring destruction of the dendrites. As the liquid percentage increases, containment of the slug becomes exceedingly difficult and therefore complicates the process. Increased amounts of energy and time are also required to heat the slugs and destroy the dendrites.\nU.S. Pat. No. 5,878,804 obviates the containment problem by heating the slug within a mold. Nevertheless, increased amounts of energy and time are still required for dendrite destruction. Cooling of the finished product is also more difficult when using a heated mold.\nFinally, U.S. Pat. No. 6,003,585 also requires the manufacture of a multialloy slug to similarly accomplish the objects of the present invention.\nTherefore, a need exists for a simplified and cost-effective thixotropic manufacturing method that can be modified to vary the properties of different parts integral to a complete assembly.\nThe present invention solves the aforementioned problems by implementing a thixotropic process under vacuum for the production of turbine rotors and other parts of intricate design that comprise high melting point alloys. The mechanical properties of semisolid forgings are tailored by microstructure or metallurgical chemistry to achieve optimized properties in specific locations of the final product.\nInitially, two or more high temperature slugs are first machined or preformed to fit within a heater in the metal forming process. The slugs are generally comprised of the same or different alloys and may be heated to a semisolid or solid state depending on design criteria.\nIn a first embodiment of the process, a first slug is heated under vacuum but retained as a solid. A second slug is also under vacuum and is heated to a thixotropic or semisolid state. Once the desired liquid/solid thixotropic ratio is attained within the second slug, the solid and semisolid slugs are forged into the mold. Upon actuation of a piston or plunger, the solid slug is driven into the semisolid slug. The semisolid slug thus enters the mold and is then exuded into predetermined areas upon the subsequent entry of the solid slug. The areas receiving the semisolid slug thereby benefit from its respective properties once the completed assembly hardens, and the areas occupied by the solid slug(s) benefit from its respective properties once the assembly is completely forged.\nThe process described is particularly well suited for manufacturing turbine assemblies comprised of integrated bore, rim, and blade components. For example, to form a rotor assembly, a first slug containing a bore alloy may be heated but remain in the solid state within a heater. A second slug, axially aligned with the first slug, contains a blade and rim alloy and is simultaneously and independently heated to a thixotropic state within the heater. After the heating step, a piston or other means drives the solid (first) slug into the thixotropic (second) slug, whereby the thixotropic slug enters the mold and is then followed by the solid slug. The solid slug thus forces the semisolid slug to exude into predetermined outer rim and blade areas of the mold. The solid slug, on the other hand, remains within the central or bore region of the finished part. It should be appreciated that one or more solid slugs may be accelerated into one or more semisolid slugs depending on design criteria.\nModified equipment design may be utilized in alternate embodiments of the high melting point semisolid process."}
{"text": "The present invention is particularly designed for use with a fall restrain harnesses for attachment to a worker in a situation where they may fall so that the shock loads from the fall can be transmitted through a structure to a suitable fixed support.\nOne problem with providing such devices is that of providing a suitable anchor at an elevated position so that the harness can be attached to save the worker from hitting the ground in a fall. The anchor point must be sufficiently elevated and sufficiently close to the worker to avoid a pendulum effect causing the falling worker to swing into contact with adjacent structures. The anchor point must accommodate the required high shock loading without damage, where the shock loading may significantly exceed any static loading. Thus the structure supporting the anchor point must have sufficient strength to meet the requirements for such high shock loading.\nIn situations where there is no suitable overhead structure, a ground or wall based system is required. In many cases there is little room or structure at the ground for the massive base required to provide the required loadings. A wall based system cannot apply high loadings to the walls without the danger of damage to the building structure.\nHowever the above system particularly designed for fall restraint, can be used for supporting other loads in a situation where excess or shock loading above predetermined level can be expected."}
{"text": "The present invention relates generally to pneumatic or hydraulic actuated relay switching systems and devices and, in a preferred embodiment thereof, more particularly provides a wholly mechanical, fluid pressure actuated sequencing relay device for regulating and controlling the order in which main and wing valve actuators are pressurized and depressurized.\nThe \"Christmas tree\" valve configuration in production fluid piping for oil and gas wells is quite common and comprises a regulated main flow supply valve and a diverting pipeline network having paths downstream from the main valve that are opened and closed by secondary valves, known as wing valves.\nSome Federal regulations now require a fail-safe valve system at the wellhead in certain oil and gas fields which is operative to assure that if the main valve is opened or closed the wing valves are automatically opened or closed, and vice versa. Therefore, it is advantageous to develop a system to perform a series of operations in a prescribed order so that a fail-safe valve system is utilized and the federal regulations are complied with.\nIt is well known in the oil and gas recovery and extraction industry that the main valve is to be opened prior to opening the wing valves to initiate production fluid flow. Similarly, when terminating the delivery of the product, it is desirable to close the wing valves first, and then close the main valve, so that the main valve is closed at a \"no flow\" condition to prolong its operating life.\nIn the past, the opening and closing sequence of the main and wing valves was effected manually. In recent years, however, the sequential operation of the main and wing valves has been automatically accomplished by a relatively complex system of mechanical sequencing circuits that incorporate a conventional block and bleed relay, a time delay mechanism, and a pilot supply for each valve actuator. It is also well known among those skilled in the art to incorporate time delay circuits between the main valve and wing valves. The circuitry initiates or delays the \"open\" and \"closed\" signal to a particular valve so as to promote the proper sequencing.\nThe major drawback with the conventional systems is an increased propensity for mechanical failure of any one of the numerous components, thereby rendering the entire system inoperable or ineffective.\nTherefore, it is desirable to develop a device to sequentially pressurize a first reservoir, before pressurizing a second reservoir, then depressurize the second before depressurizing the first, thereby eliminating the need for the multiplicity of elements in the conventional relay switching system. It is accordingly an object of the present invention to provide such a device for use in a relay switching system."}
{"text": "Fire alarms and fire alarm systems are generally known. Such systems generally include a number of fire detectors distributed around a protected area. The fire detectors may be based upon any of a number of different fire detection technology (e.g., smoke detection, carbon monoxide detection, etc.).\nEach of the fire detectors is typically connected to a fire alarm control panel. The connection between each of the sensors and the control panel may be wired or wireless.\nThe control panel monitors each of the sensors for an indication of the presence of a fire and, in response, sounds an alarm. The control panel may also send notification of the fire to a central monitoring station via a communication connection (e.g., a dial-up connection, the Internet, etc.).\nMost fire alarm control panels are typically provided with a display that provides an indication of any sensors that have been activated by a fire. The indications are typically provided with an alpha-numeric identifier or a short description of the location of the fire.\nWhile such systems are effective for personnel familiar with the protected facility, they are not very helpful for outside firefighters. In this case, outside firefighters may require access to a facility map to a cross reference between the identifier of a fire sensor to a location within the protected facility.\nHowever, even with the activated sensor identified on a map, the firefighter may still not be able to quickly access the fire. Doors may be locked. Corridors may be blocked. Accordingly, a need exists for better methods of guiding firefighters to fires."}
{"text": "Many diseases of the central nervous system are influenced by the adrenergic, the dopaminergic, and the serotonergic neurotransmitter systems. For example, serotonin has been implicated in a number of diseases and conditions which originate in the central nervous system. A number of pharmacological and genetic experiments involving receptors for serotonin strongly implicate the 5-HT2c receptor subtype in the regulation of food intake (Obes. Res. 1995, 3, Suppl. 4, 449S-462S). The 5-HT2c receptor subtype is transcribed and expressed in hypothalamic structures associated with appetite regulation. It has been demonstrated that the non-specific 5-HT2c receptor agonist m-chlorophenylpiperazine (mCPP), which has some preference for the 5-HT2c receptor, causes weight loss in mice that express the normal 5-HT2c receptor while the compound lacks activity in mice expressing the mutated inactive form of the 5-HT2c receptor (Nature 1995, 374, 542-546). In a recent clinical study, a slight but sustained reduction in body weight was obtained after 2 weeks of treatment with mCPP in obese subjects (Psychopharmacology 1997, 133, 309-312). Weight reduction has also been reported from clinical studies with other “serotonergic” agents (see e.g. IDrugs 1998, 1, 456-470). For example, the 5-HT reuptake inhibitor fluoxetine and the 5-HT releasing agent/reuptake inhibitor dexfenfluramine have exhibited weight reduction in controlled studies. However, currently available drugs that increase serotonergic transmission appear to have only a moderate and, in some cases, transient effects on the body weight.\nThe 5-HT2c receptor subtype has also been suggested to be involved in CNS disorders such as depression and anxiety (Exp. Opin. Invest. Drugs 1998, 7, 1587-1599; IDrugs, 1999, 2, 109-120).\nThe 5-HT2c receptor subtype has further been suggested to be involved in urinary disorders such as urinary incontinence (IDrugs, 1999, 2, 109-120).\nCompounds which have a selective effect on the 5-HT2c receptor may therefore have a therapeutic potential in the treatment or prophylaxis of disorders like those mentioned above. Of course, selectivity also reduces the potential for adverse effects mediated by other serotonin receptors.\nExamples of such compounds are (2R)-1-(3-{2-[(2-ethoxy -3-pyridinyl)oxy]ethoxy}-2-pyrazinyl)-2-methylpiperazine, (2R)-methyl-1-{3-[2-(3-pyridinyloxy]ethoxy]-2-pyrazinyl}piperazine and pharmaceutically acceptable acid salts thereof. WO 00/76984 (hereinafter called D1) relates to a process for the preparation of such compounds on a small scale such as a gram scale. A problem to be solved by the present invention was to prepare such compounds on a large scale such as on a kilogram scale. The following factors are more important for preparation on a large scale, in comparison to preparation on a small scale:                to obtain a high yield of the desired products for economy reasons,        that the processes for preparation are safe with regard to explosion,        that the reagents and solvents used are relatively non-toxic,        that the products obtained are relatively stable, and        that the reaction times are relatively short.        \nThese problems have been solved by the present invention. It has been shown that the yields of the desired products according to the present invention are higher than the yields according to D1. In the experimental part, the yields according to the present invention and D1 are compared. Regarding the choice of solvents for the process steps, dioxane, as used according to D1, has been replaced by solvents such as MtBE and THF (see step (ii) below) which are less prone to form peroxides and which are less carcinogenic than dioxane. Furthermore, it has been shown that (2R)-methyl-1-{3-[2-(3-pyridinyloxy)ethoxy]-2-pyrazinyl}piperazine, L-malate salt prepared according to the present invention (see Example 3A) has superior properties compared to (2R)-methyl-1-{3-[2-(3-pyridinyloxy)ethoxy]-2-pyrazinyl}piperazine, hydrochloride prepared according to D1 in that the former has a higher crystallinity, is less hygroscopic and has a higher chemical stability than the latter. Regarding chemical stability, D1 discloses the preparation of (2R)-1-(3-chloro-2-pyrazinyl) -2-methylpiperazine, which has also been prepared in Example 8 below. This compound is not stable for storage as a free base. As (2R)-1-(3-chloro-2-pyrazinyl)-2-methylpiperazine is a key intermediate, chemical stability during long term storage is important with regard to process economy. It has now been found that the corresponding hydrochloride salt thereof is considerably more stable, which has been prepared in Example 9 below.\nThe method to prepare Example 3C is a good way of increasing the purity of Example 2C. It has been shown that Example 2C with a purity of 60-70% gives Example 3C with a purity of 99% in one crystallization step. By contrast, the same purity increasing effect has not been achieved by making the acetate of Example 2C.\nRegarding reaction time, the reaction according to Example 2A below was complete at room temperature in 15 minutes. The same compound has been prepared in Example 173 in D1. The procedure of Example 172, step 3 has been followed, wherein the reaction was stirred at 85° C. for 15 h. Furthermore, the reaction according to Example 2B below was complete in 35 minutes at 55° C. The same compound has been prepared in Example 200 in D1. The procedure of Example 192, step 3 has been followed, wherein the reaction was stirred at 90° C. for 2 h."}
{"text": "Adhesives are widely used for labeling and packaging, in laminating objects such as paper, plastic, wood, and metal, as well as in electronics manufacturing where specialty adhesives are applied for the assembly of electronic components. For example, adhesives are used to create permanent or temporary graphic overlays, security documents, and decals. These adhesives are typically pre-applied on the desired parts using conventional coating or screen printing techniques, which are associated with high materials usage, high tolling and handling costs, and long production cycle time. In addition, the screen printable adhesives contain solvent media or low molecular weight monomers, which require drying or a curing step. Therefore, there is a need for a method for applying adhesives, which allows disposing the adhesive selectively on the surface on-demand, where the bonding is required, using digital printing techniques."}
{"text": "Compilers are generally used to transform one representation of a computer program into another representation. Typically, but not exclusively, compilers are used to transform a human readable form of a program such as source code into a machine readable form such as object code.\nA computer program suitable for compilation by a compiler is composed of a series of \"statements\". Some statements generate, modify, retrieve or store information. Other statements may control the flow of the program, for example, by testing the value of a variable and causing program flow to continue in different directions based on the value of the variable. In most programs of any significant length, the statements are collected into \"procedures\", which perform well-defined functions and can be used in potentially multiple places with the program. Frequently, the procedures in a large program are further collected into \"modules\", each of which is responsible for a particular major subset of the functions of the program. In a program structure of this kind, the compiler is used to compile the modules individually, after which the compiled modules are \"linked\" together to form a single, cohesive computer program. This approach allows the programmer to upgrade or debug, and then recompile, each module separately, without need for recompiling the other modules.\nOne type of compiler is an optimizing compiler which includes an optimizer for enhancing the performance of the machine readable representation of a program. Some optimizing compilers are separate from a primary compiler, while others are built into a primary compiler to form a multi-pass compiler. Both types of compilers may operate either on a human readable form, a machine readable form, or any intermediate representation between these forms.\nOne optimization technique is known as \"profiling\" the program. A program is profiled by compiling the program, and delivering it to a test environment which simulates actual field operation of the program. While the program operates in the test environment, records are kept on the extent to which certain sections of the program are used. After the test has been completed, the profile records are used by an optimizing compiler, to recompile the program in a manner which enhances the efficiency of the program. For example, one known technique is to place sections of the program which are used at approximately the same time, in nearby memory locations, so as to speed access to the program.\nA common computer programming approach is known as procedural programming. In procedural programming, a program is broken into many small procedures, each including a sequence of statements (and in some cases, data), and each of which is responsible for particular well-defined activities. The procedures are invoked when particular actions are needed. Typically, procedures can invoke each other, as part of operation of the program. In such a situation, the procedure which is invoked is typically referred to as the \"child\" procedure, and the procedure which invokes the child procedure is referred to as the \"parent\" procedure.\nWhile procedural programming can simplify programming effort and reduce complexity, one of the unfortunate results of a highly procedural computer program, is that the program, when operating, frequently transfers control between the various procedures (executes \"procedure calls\"). This creates a substantial overhead, in that each transfer of control between procedures requires multiple computer operations, both to transfer flow control to a procedure and to return flow control from the procedure.\nA similar unfortunate result occurs in so-called \"object oriented\" programming. In object oriented programming, data and a set of procedures (called \"methods\") are encapsulated together, and only the procedures encapsulated with data are permitted to modify that data. This style of programming naturally causes procedure calls to proliferate, typically to a greater extent than procedural programming.\nTo address the problem of high procedure call overhead, modern compilers optimize programs so as to avoid procedure calls. One optimization approach is to \"inline\" procedures, that is, to copy the entire body of the child procedure, into the body of the parent procedure, at each location in the parent procedure where the child procedure is referenced. This is usually done only when the child procedure is relatively small and is called from relatively few locations, in order to minimize the extent to which the overall compiled program size is increased due to inlining."}
{"text": "The financial and legal world is quite complex. Cross-border financial transactions among entities are commonplace, resulting in entities having multiple financial exposures with others around the world. This can be difficult to manage and can lead to many problems. For example, a financial institution may have large intra-day foreign exchange settlement obligations with a number of different trading partners. A large financial institution may have millions of dollars of exposure to their largest counterparties on any given day. Entities may also have large exposures based on counterparty credit risk and liquidity risk. Many entities enter into agreements to manage and control these exposures.\nWhen trading partners agree to offset their positions or obligations, they are “netting”. By doing so, they reduce a large number of individual positions or obligations to a smaller number of positions or obligations, and it is on this netted position that the two trading partners settle their outstanding obligations. Besides reducing transaction costs and communication expenses, netting is important because it reduces credit and liquidity risks, and ultimately systemic risk.\nNetting agreements have been used to manage these exposures in a number of different contexts. Netting agreements are a contractual mechanism to offset payables against receivables to reduce an entity's exposure to a counterparty. Netting agreements are used, for example, to reduce credit exposure to the net obligation of a counterparty. The enforceability and use of netting agreements varies by jurisdiction. For example, in the United States, netting in bankruptcy or insolvency is enforceable under the federal bankruptcy code. Netting between United States-based counterparties is permitted by the Financial Institutions Reform, Recovery, and Enforcement Act of 1989 (FIRREA). Different jurisdictions have different rules and laws regarding the use and enforceability of netting agreements. This can make it quite difficult for an entity to manage credit and exchange risk with any certainty. Frequently, an opinion of counsel regarding the legality of a particular netting relationship is required for each jurisdiction, and often for each netting agreement.\nIt would be desirable to provide a system which allows the analysis of netting agreements. It would further be desirable to provide a system which allows the automated analysis of netting agreements. It would further be desirable to provide a system which allows a number of issues associated with agreements to be analyzed based on the particular fact pattern of each agreement. It would further be desirable to provide a system which allows the analysis of agreements and which allows updating of netting positions between parties."}
{"text": "The invention described herein may be manufactured and used by or for the Government of the United States of America for governmental purposes without payment of any royalties thereon or therefor.\n1. Field of the Invention\nThe present invention relates to interface gun mounts and, more particularly, to a lightweight machine gun and ammunition can mount for use with weapons that do not have a forward mount point or a mid weapon mount orientation is required.\n2. Description of the Background\nThe size and weight of many firearms, particularly large guns, such as machine guns, precludes accuracy and stability without some type of support or mounting apparatus to hold the weapon steady while being fired. Thus, mounting devices and other support apparatus have long been used with large guns to stabilize the gun and reduce vibration, thereby improving accuracy and alleviating fatigue and discomfort of the shooter. Indeed, when mounting devices and other support apparatus have not been available a shooter has often needed to rely on any immovable object available, such as walls, rocks, tree trunks, etc.\nA variety of interface elements have evolved for use in mounting rifles or light machine guns. For example, U.S. Pat. No. 1,273,178 to Heinemann shows an apparatus designed for attachment to a hand machine gun to counteract its tendency to jump/recoil during firing. The apparatus is acted upon in a downward direction by the gases leaving the muzzle, thereby creating a downward force at the forward end of the gun\"\"s barrel. The apparatus also includes a 10 pair of downward extending eyelets for optional attachment to a gun rest or mount.\nU.S. Pat. No. 1,273,178 to Perry et al. discloses an apparatus utilized to mount a rapid fire gun on a motorcycle. The mounting apparatus allows a single individual to both drive the motorcycle and fire the gun. The mounting of the gun is such that it may be easily operated and adjusted without the rider leaving the seat of the motorcycle. The apparatus includes a circular clamping attachment to fixedly attach the gun to the position adjustment mechanism.\nU.S. Pat. No. 5,194,678 to Kramer discloses a rest for the firearm or the like that attaches to the sling on the forearm of the firearm. The firearm rest is comprised of two major assemblies. The first assembly attaches to the sling swivel and provides a surrogate sling swivel and a female receptacle for the second assembly. The second assembly consists of a male protrusion that mates with the first assembly and two lightweight legs that can be quickly assembled and disassembled.\nU.S. Pat. No. 5,711,103 to Keng discloses a bipod mounting assembly for attaching a bipod to the forearm stock portion of a firearm. The apparatus includes a mounting yoke adapted to quickly and easily attach to the swivel stud connector mounted to the forearm stock portion of the firearm. The mounting yoke is adapted to receive a mounting block thereover, with the mounting block being attached to the mounting yoke to thus attach the mounting block to the forearm stock portion of the firearm. A bipod-mounting frame is releasably attachable to the mounting block by a quick-release locking catch to enable the quick attachment/detachment of the legs of the bipod from the mounting block, and thus the firearm.\nAlthough each of the prior art examples provides an interface for use in mounting firearms which uses the pre-existing structure of the gun, they do not accommodate weapons that have no forward attach point and which are not designed for mounting.\nConsequently, there is a significant need for a lightweight machine gun and ammunition can gun mount for use with weapons that do not have a forward mount point or a mid weapon mount orientation is required, such as the MK46 and M249SPW weapons.\nFIG. I is a drawing of an existing MK46 machine gun which illustrates the location of the ammunition can which is attached to an existing ammunition can mount bracket of the weapon, and a supplied forearm rail system. In addition, the weapon includes a main mounting lug located proximate to the trigger guard (obscured) for a single-point mount. Otherwise, no provision for vehicle or tripod mounting is provided.\nThe design challenge with such guns is to determine how to securely attach a vehicle or tripod mount to the weapon to allow easy firing without impacting the performance of the gun.\nIt is, therefore, the object of the present invention to provide a machine gun and ammunition can interface gun mount for use with weapons that do not have a forward mount point or a mid weapon mount orientation is required.\nIt is yet another object to provide a machine gun and ammunition can interface gun mount with two point secure attachment to allow easy firing without impacting the performance of the gun.\nIt is another object to provide a machine gun and ammunition can interface gun mount that is small, lightweight and economical to produce.\nIt is still another object to provide a machine gun and ammunition can interface gun mount with two point attachment for improved stability and anti-rotation, the first point being the ammunition can mount attached to the weapon for anti-rotation and the second point being directly in front of the trigger mechanism.\nIt is yet another object to provide a machine gun and ammunition can interface gun mount in which the ammunition can is repositioned to the left side of the gun for better alignment and ammunition feed.\nAccording to the present invention, the above-described and other objects are accomplished by a gun mount that is securely attached to the weapon by using two existing separate structural features of the gun as attach points. The first point of attachment is the ammunition can mounting bracket and the second attachment point is the main mounting lug located proximate the trigger guard. The two point attachment allows the shooter to steady the gun and eliminate vibration and improve accuracy of aim. The ammunition can is relocated from the common position underneath the gun to a vertical position on the left side of the gun for better alignment and ammunition feed. Both the ammunition can interface and main mounting lug interface are formed of coated aluminum to effect a strong, compact and lightweight design. Thus, the present invention is small, compact, lightweight, and economical to manufacture."}
{"text": "1. Technical Field\nThis invention relates generally to the field of video surveillance services. More specifically, this invention relates to video surveillance services for video image processing, organization, and archival.\n2. Description of the Related Art\nVideo surveillance technology is an industry that has become available to everyday end users. The wide variety of video surveillance tools may span from webcams for home surveillance to full service business video surveillance and security systems with video security cameras. As well, the video surveillance industry includes algorithms for feature detection on images and storage of video, such as on YouTube by YouTube, LLC in San Bruno, Calif. The industry includes Wi-Fi video cameras for deployment in the home, such as for example, by Dropcam, Inc. in San Francisco, Calif. By such camera, an end user may monitor a particular scene, such as a room in a house, from a handheld device or a computer. Another current product offered in the industry is a competing video surveillance archiving service by sensr.net, Inc. in Incline Village, Nev. By such service, an end user may monitor a scene from an online account and receive alerts activated by detected motion. Further, images from such camera may be stored in the cloud."}
{"text": "Streaming media generally refers to media content that is, or at least may be, played via playback software or a playback device at the same time that the media content is being downloaded from a source such as a media server. Streaming media content, e.g., video and audio content, may be provided according to a variety of standards and formats. For example, video standards such as QuickTime and RealMedia, and also standards promulgated by the Motion Picture Experts Group (MPEG), etc. are well known.\nMany standards for streaming media content, such as MPEG streaming content delivery, were designed with dual objectives of (1) preserving network bandwidth and (2) maintaining video quality. However, MPEG and most other kinds of media streams are not designed with an objective of timely delivery of content, e.g., diminishing latency. For example, at present, content processing devices such as set top boxes (STBs) generally use MPEG and are designed to deliver a complete video stream at the expense of latency.\nA media stream such as an MPEG stream reaching a STB can potentially face network jitter which can cause excessive storage of media frames, e.g., video frames, in a buffer in the STB. Frames received by the STB are not displayed until all the prior frames are displayed. However, in some contexts, e.g., interactive applications such as gaming, users expect timely updates on their video displays in response to a key press. At present, latencies caused by jitter often leave users with a video stream that is unsatisfactory for supporting applications in a variety of contexts.\nOnline gaming is one context in which media stream latencies may result in an unsatisfactory user experience. For example, when gaming is provided through a content processing device such as a set top box (STB) or the like, a game session is delivered as an MPEG video stream or the like through a packet switched network from a game server in a Video Hub Office (VHO) to an STB in a customer premises. That is, the game session is conducted on the game server but is presented, through the MPEG stream, by the STB. Accordingly, the game session is encoded as an MPEG stream and streamed to the STB over the network. The MPEG stream is decoded by the STB and then displayed on a media playback device such as a television or video monitor. User inputs to the game are gathered through an input device such as a radio frequency (RF) or infrared remote control, a universal serial bus (USB) gamepad, etc. User inputs are then sent back to the game server over the packet switched network. The game server receives the user inputs and provides them to the game session for processing, thereby altering the output video stream where appropriate based on the inputs. Latencies in the MPEG stream may cause user inputs to be ill timed and/or ineffective, thus rendering the gaming experience unsatisfactory for the user.\nThus, many standards for providing streaming media, such as MPEG streaming content delivery, were designed with traditional objectives of (a) preserving network bandwidth and (2) maintaining video quality. However, as is the case with many media streams, MPEG video streams are not designed with an objective of timely delivery of content, e.g., diminishing latency. For example, at present, content processing devices such as set top boxes (STBs) are designed to deliver a complete video stream at the expense of latency."}
{"text": "Iron-based particles have long been used as a base material in the manufacture of structural components by powder metallurgical methods. The iron-based particles are first molded in a die under high pressures in order to produce the desired shape. After the molding step, the structural component usually undergoes a sintering step to impart the necessary strength to the component.\nMagnetic core components have also been manufactured by such power metallurgical methods, but the iron-based particles used in these methods are generally coated with a circumferential layer of insulating material.\nTwo important characteristics of an iron core component are its magnetic permeability and core loss characteristics. The magnetic permeability of a material is an indication of its ability to become magnetized, or its ability to carry a magnetic flux. Permeability is defined as the ratio of the induced magnetic flux to the magnetizing force or field intensity. When a magnetic material is exposed to a rapidly varying field, the total energy of the core is reduced by the occurrence of hysteresis losses and/or eddy current losses. The hysteresis loss is brought about by the necessary expenditure of energy to overcome the retained magnetic forces within the iron core component. The eddy current loss is brought about by the production of electric currents in the iron core component due to the changing flux caused by alternating current conditions.\nEarly magnetic core components were made from laminated sheet steel, but these components were difficult to manufacture and experienced large core losses at higher frequencies. Application of these lamination-based cores is also limited by the necessity to carry magnetic flux only in the plane of the sheet in order to avoid excessive eddy current losses. Sintered metal powders have been used to replace the laminated steel as the material for the magnetic core component, but these sintered parts also have high core losses and are restricted primarily to direct current operations.\nResearch in the powder metallurgical manufacture of magnetic core components using coated iron-based powders has been directed to the development of iron powder compositions that enhance certain physical and magnetic properties without detrimentally affecting other properties. Desired properties include a high permeability through an extended frequency range, high pressed strength, low core losses, and suitability for compression molding techniques.\nWhen molding a core component for AC power applications, it is generally required that the iron particles have an electrically insulating coating to decrease core losses. The use of a plastic coating over the iron particles (see U.S. Pat. No. 3,935,340 to Yamaguchi) and the use of doubly-coated iron particles (see U.S. Pat. No. 4,601,765 to Soileau et al.) have been employed to insulate the iron particles and therefore reduce eddy current losses.\nRecently, it has been found that the insulating polymeric material need not be present in the powder composition as a full coating of the individual iron particles, but rather can be present in the form of discrete particles that are integrally admixed with the iron particles. The present invention is directed to a method of forming this admixture in a manner that ensures homogeneity and thereby leads to improved magnetic properties of pressed parts made with the powder composition. The invention eliminates the need to provide the iron-based particles with a circumferential coating of the polymeric material, which coating generally required the use of more expensive fluidized bed processes."}
{"text": "Montelukast sodium is the active pharmaceutical ingredient of SINGULAIR®, and is approved for the treatment of asthma and allergic rhinitis. The molecular structure of montelukast is as shown below:\n\nMontelukast sodium is described in U.S. Pat. No. 5,565,473. A crystalline form of montelukast sodium (hereinafter referred to as “Form A”) is described in U.S. Pat. No. 5,614,632."}
{"text": "This invention relates generally to hygienic tools and more particularly to tools used to clean the hands and fingers.\nNot only does civilized society demand cleanliness, but, basic hygiene also requires regular and thorough cleaning of the body. While the majority of the skin is either easily cleaned or is not excessively exposed to dirt and disease, the same cannot be said of the hands.\nBy their very nature, the hands are exposed in daily life and work to a wide range of grime and disease. The fingers pose the most difficult scenario for cleaning. The fingers provide extensive surface space which is not readily accessed with a hand brush; and the fingernails are particularly difficult to clean. Even after a thorough washing, a mechanic often still has xe2x80x9cblackxe2x80x9d fingernails.\nWhether the cleaning required is after work in the garden or on the car, or requires antiseptic levels for medical procedures, cleaning of the fingers and fingernails is particularly difficult and time consuming. Since the task is so very time consuming, often the cleaning is haphazardly done.\nIt is clear from the foregoing that there is a need for an efficient tool to assist in the cleaning of fingers and fingernails.\nThe invention is a tool for cleaning fingers. In one embodiment of the invention, a closed-ended tube is used. The tube is large enough to accept a to-be-cleaned finger into the tube. Ideally, the tube\"\"s length is not longer than the shortest of the user\"\"s fingers; this attribute will become clear later.\nThe preferred embodiment has two interlocking tubes. Both tubes are closed ended with the open ends mating/connecting with each other. This interlocking attribute provides for an easily xe2x80x9csealedxe2x80x9d or contained apparatus for transport or storage between uses.\nEach tube is configured to accept a finger from the user. Ideally, the first tube is used by the user initially with the second tube performing an optional cleaning or buffing function.\nWithin the first tube is a longitudinal brush. The longitudinal brush is used to clean the sides of the inserted finger as the tube is rotated around the finger. The longitudinal brush is secured to the tube through a variety of ways, obvious to those of ordinary skill in the art.\nTwo such methods include: securing the longitudinal brush to the interior wall of the tube; and, securing an end of the longitudinal brush to the closed end of the tube.\nIn the latter situation, where the longitudinal brush is secured to the end of the tube, in some embodiments the end of the tube is removable. The ability to remove the end of the tube (and hence the longitudinal brush) permits the brush and tube to be easily cleaned, or the replacement of the longitudinal brush.\nIn the preferred embodiment, an end brush is secured to the interior portion of the closed end of the tube. This end brush permits the user to press the fingernail against the end brush. When the tube is rotated, the end brush moves against the surface and end of the fingernail to clean the fingernail. This action provides excellent cleaning beneath the fingernail.\nIn the embodiment where two tubes are used, the second tube includes a soft brush is used to softly clean the finger. This tube xe2x80x9cbuffsxe2x80x9d the finger and fingernail to provide a finishing cleaning operation.\nHoles in the tube permit soap to enter and be used for the cleaning process. When the tube is in use, the user places the tube and finger under the stream of water and soaps the assembly, allowing the soap and water to enter the tube to assist in the cleaning process as the tube is rotated.\nIn one embodiment, a cavity is created within the tube. This cavity is xe2x80x9cchargedxe2x80x9d or loaded with soap, lotion, or antibacterial agents, using an exterior opening or portal. The liquid within the cavity is dispensed from the cavity into the interior of the tube; then, as the tube is rotated, the soap, lotion, or antibacterial lotion heightens the cleaning process.\nTo provide an antiseptic environment, some embodiments of the invention provide for a coating of anti-bacterial agents on the bristles on the brushes. This embodiment provides an anti-bacterial agent without any effort on the part of the user; periodically, the entire unit is replaced or the brushes are replaced (as in the case of the longitudinal brushes connected to a removable end-cap).\nThe invention, together with various embodiments thereof, will be more fully explained by the accompanying drawings and following description thereof."}
{"text": "The instant invention relates generally to outdoor watering devices and more specifically it relates to a lawn sprinkler for a garden hose.\nNumerous outdoor watering devices have been provided in prior art that are adapted to spray water through rotating nozzles. For example U.S. Pat. Nos. 172,024; 270,664 and 1,558,355 all are illustrative of such prior art. While these units may be suitable for the particular purpose to which they address, they would not be as suitable for the purposes of the present invention as heretofore described."}
{"text": "The present invention is directed to polishing pads used for creating a smooth, ultra-flat surface on such items as glass, semiconductors, dielectric/metal composites, magnetic mass storage media and integrated circuits. More specifically, the present invention relates to grafting and preserving the grafted surface of polymers, preferably thermoplastic foam polymers, thereby transforming their mechanical and chemical properties to create more suitable polishing pads therefrom.\nChemical-mechanical polishing (CMP)is used extensively as a planarizing technique in the manufacture of VLSI integrated circuits. It has potential for planarizing a variety of materials in IC processing but is used most widely for planarizing metallizied layers and interlevel dielectrics on semiconductor wafers, and for planarizing substrates for shallow trench isolation.\nIn trench isolation, for example, large areas of field oxide must be polished to produce a planar starting wafer. Integrated circuits that operate with low voltages, i.e., 5 volts or less, and with shallow junctions, can be isolated effectively with relatively shallow trenches, i.e., less than a micron. In shallow trench isolation (STI) technology, the trench is backfilled with oxide and the wafer is planarized using CMP. The result is a more planar structure than typically obtained using LOCOS, and the deeper trench (as compared with LOCOS) provides superior latch up immunity. Also, by comparison with LOCOS, STI substrates have a much reduced xe2x80x9cbirds\"\" beakxe2x80x9d effect and thus theoretically provide higher packing density for circuit elements on the chips. The drawbacks in STI technology to date relate mostly to the planarizing process. Achieving acceptable planarization across the full diameter of a wafer using traditional etching processes has been largely unsuccessful. By using CMP, where the wafer is polished using a mechanical polishing wheel and a slurry of chemical etchant, unwanted oxide material is removed with a high degree of planarity.\nSimilarly, integrated circuit fabrication on semiconductor wafers require the formation of precisely controlled apertures, such as contact openings or xe2x80x9cvias,xe2x80x9d that are subsequently filled and interconnected to create components and very large scale integration (VLSI) or ultra large scale integration (ULSI) circuits. Equally well known is that the patterns defining such openings are typically created by optical lithographic processes that require precise alignment with prior levels to accurately contact the active devices located in those prior levels. In multilevel metallization processes, each level in the multilevel structure contributes to irregular topography. In three or four level metal processes, the topography can be especially severe and complex. The expedient of planarizing the interlevel dielectric layers, as the process proceeds, is now favored in many state of the art IC processes. Planarity in the metal layers is a common objective, and is promoted by using plug interlevel connections. A preferred approach to plug formation is to blanket deposit a thick metal layer on the interlevel dielectric and into the interlevel windows, and then remove the excess using CMP. In a typical case, CMP is used for polishing an oxide, such as SiO2, Ta2O5, W2O5. It can also be used to polish nitrides such as Si3N4, TaN, TiN, and conductor materials used for interlevel plugs, such as W, Ti, TIN.\nCMP generally consists of the controlled wearing of a rough surface to produce a smooth specular finished surface. This is commonly accomplished by rubbing a pad against the surface of the article, or workpiece, to be polished in a repetitive, regular motion while a slurry containing a suspension of fine particles is present at the interface between the polishing pad and the workpiece. Commonly employed pads are made from felted or woven natural fibers such as wool, urethane-impregnated felted polyester or various types of filled polyurethane plastic.\nA CMP pad ideally is flat, uniform across its entire surface, resistant to the chemical nature of the slurry and have the right combination of stiffness and compressibility to minimize effects like dishing and erosion. In particular, there is a direct correlation between lowering Von Mises stress distributions in the pad and improving polishing pad removal rates and uniformity. In turn, Von Mises stresses may be reduced though the controlled production of pad materials of uniform constitution, as governed by the chemical-mechanical properties of the pad material.\nCMP pad performance optimization has traditionally involved the empirical selection of materials and use of macro fabrication technologies. For example, a pad possessing preexisting desirable porosity or surface texture properties may be able to absorb particulate matter such as silica or other abrasive materials. Or, patterns of flow channels cut into the surface of polishing pads may improve slurry flow across the workpiece surface. The reduction in the contact surface area effected by patterning also provides higher contact pressures during polishing, further enhancing the polishing rate.\nAlternatively, intrinsic microtextures may be introduced into pads by using composite or multilayer materials possessing favorable surface textures as byproduct of their method of manufacture. Favorable surface microtextures may also be present by virtue of bulk non-uniformities introduced during the manufacturing process. When cross-sectioned, abraded, or otherwise exposed, these bulk non-uniformities become favorable surface microtextures. Such inherent microtextures, present prior to use, may permit the absorption and transport of slurry particles, thereby providing enhanced polishing activity without the need to further add micro- or macrotextures.\nThere are, however, several deficiencies in polishing pad materials selected or produced according to the above-described empirical techniques. Pads made of layers of polymer material may have thermal insulating properties, and therefore unable conduct heat away from the polishing surface, resulting in undesirable heating during polishing. Numerous virgin homogenous sheets of polymers such as polyurethane, polycarbonate, nylon, polyureas, felt, or polyester, have poor inherent polishing ability, and hence not used as polishing pads. In certain instances, mechanical or chemical texturing may transform these materials, thereby rendering them useful in polishing.\nHowever, polyurethane based pads, currently in widespread use, are decomposed by the chemically aggressive processing slurries by virtue of the inherent chemical nature of urethane. This decomposition produces a surface modification in and of itself in the case of the polyurethane pads.\nYet another approach involves modifying the surface of CMP polishing pads materials to improve the wetability of the pad surface, the adhesion of surface coatings, and the application performance of these materials. Plasma treatment of polishing pad materials is one means to functionalize and thereby modify polishing pad surfaces. However, the simple functionalization of pad surfaces by plasma treatment is known to be a temporary effect, with spontaneous loss of functionalization after one to two days. While some success in the preservation of functionalized pad surfaces has been obtained for fluorinated polymeric surfaces, this has not been demonstrated for other polymeric surfaces, and in particular, thermoplastics.\nAccordingly, what is needed in the art is an improved process for functionalizing and preserving a semiconductor wafer thermoplastic polishing pad surface, thereby providing a rapid rate of polishing and yet reducing scratches and resultant yield loss during chemical/mechanical planarization.\nTo address the deficiencies of the prior art, the present invention, in one embodiment, provides a polymer, preferably thermoplastic foam polymer, comprising a thermoplastic foam substrate having a modified surface thereon and a grafted surface on the modified surface.\nIn another embodiment, the present invention provides a method for preparing a polymer, preferably a thermoplastic foam polymer. The method comprises the steps of providing a thermoplastic foam substrate, exposing the substrate to an initial plasma reactant to produce a modified surface thereon, and exposing the modified surface to a secondary plasma reactant to create a grafted surface on the modified surface.\nYet another embodiment provides a method of manufacturing a polishing pad. The method comprises providing a thermoplastic foam substrate, and then forming a thermoplastic foam polishing body with a grafted surface by including those steps described above. A polishing pad is then formed from the thermoplastic foam polishing body that is suitable for polishing a semiconductor wafer or integrated circuit using the grafted surface.\nIn still another embodiment, the present invention provides a polishing apparatus. This particular embodiment includes a mechanically driven carrier head, a polishing platen, and a polishing pad attached to the polishing platen. The carrier head is positionable against the polishing platen to impart a polishing force against the polishing platen. The polishing pad includes a polishing body comprising a material wherein the material is a thermoplastic foam polymer.\nThe foregoing has outlined, rather broadly, preferred and alternative features of the present invention so that those skilled in the art may better understand the detailed description of the invention that follows. Additional features of the invention will be described hereinafter that form the subject of the claims of the invention. Those skilled in the art should appreciate that they can readily use the disclosed conception and specific embodiment as a basis for designing or modifying other structures for carrying out the same purposes of the present invention. Those skilled in the art should also realize that such equivalent constructions do not depart from the spirit and scope of the invention in its broadest form."}
{"text": "1. Field of the Invention\nThe present invention relates to a seeding machine and in particular to a lifting mechanism for a row cleaner that lifts the row cleaner upon lifting of an associated furrow opener.\n2. Description of the Prior Art\nU.S. Pat. No. 5,341,754 discloses a grain drill assembly with several grain drills spaced along a transverse support. Each grain drill contains a furrow opener that forms a furrow for the deposit of seed. Each furrow opener is attached to the transverse support to pivot vertically. A row cleaner is located ahead of the furrow opener and is mounted to the transverse support through a parallelogram linkage to pivot vertically. The row cleaner is forced against the ground by means of a spring.\nU.S. Pat. No. 5,878,678 discloses that the row cleaner be attached to the furrow opener so that the latter carries the row cleaner when it is raised. While it is advantageous to lift the row cleaner when the furrow opener is raised, with the arrangement shown in the '678 patent, the row cleaner is not raised to the same degree as the furrow opener. This is particularly disadvantageous during operation on the road. In addition, the row cleaner always engages the ground when the furrow opener is in the lowered position. There is no ability to operate the furrow opener without using the row cleaner. When seeding in a prepared seedbed, the row cleaner is not needed."}
{"text": "The present disclosure relates generally to the field of orthopedics and spinal surgery, and in some embodiments, the present disclosure relates to intervertebral prosthetic joints for use in the total or partial replacement of a natural intervertebral disc, and methods and tools for use therewith.\nIn the treatment of diseases, injuries or malformations affecting spinal motion segments, and especially those affecting disc tissue, it has long been known to remove some or all of a degenerated, ruptured or otherwise failing disc. In cases involving intervertebral disc tissue that has been removed or is otherwise absent from a spinal motion segment, corrective measures are taken to ensure the proper spacing of the vertebrae formerly separated by the removed disc tissue.\nIn some instances, the two adjacent vertebrae are fused together using transplanted bone tissue, an artificial fusion component, or other compositions or devices. Spinal fusion procedures, however, have raised concerns in the medical community that the bio-mechanical rigidity of intervertebral fusion may predispose neighboring spinal motion segments to rapid deterioration. More specifically, unlike a natural intervertebral disc, spinal fusion prevents the fused vertebrae from pivoting and rotating with respect to one another. Such lack of mobility tends to increase stresses on adjacent spinal motion segments.\nAdditionally, several conditions may develop within adjacent spinal motion segments, including disc degeneration, disc herniation, instability, spinal stenosis, spondylosis and facet joint arthritis. Consequently, many patients may require additional disc removal and/or another type of surgical procedure as a result of spinal fusion. Alternatives to spinal fusion are therefore desirable.\nIn particular, this disclosure relates to an articulating disc prosthesis that can be inserted from the anterior approach to aid in the correction of spondylolisthesis."}
{"text": "1. Field of the Invention\nThis invention relates to a panel for supporting a plurality of cards in a predetermined, orderly manner and with the supported cards being anchored relative to the panel against shifting relative thereto due to wind, vibration or upset of the panel. The panel may be constructed to support only the cards being played by a single player or all of the played cards during a game by a plurality of game players.\n2. Description of Related Art\nVarious different forms of card holding game panels, card game boards, game packages and card game playing aids heretofore have been provided including some of the general structural and operational features of the instant invention. Examples of these different devices are disclosed in U.S. Pat. Nos. 1,880,175, 2,115,276, 2,150,850, 2,450,325, 3,667,757, 4,317,515 and 4,436,324. However, these previously known devices do not include the overall combination of structural and operational features of the instant invention which coact to provide a simple, readily constructed, durable and easy to use playing card holder."}
{"text": "The present invention generally relates to a data reproducing circuit for a memory system having a reading head, such as a magnetic disk memory system, a magnetic tape memory system, an optical memory system, etc., and reproducing analog signals read out from the head as digital reproduction signals, with a high accuracy.\nA conventional data reproducing circuit for a magnetic disc device is designed to detect the position of a peak of a reproduction signal from a magnetic disk through a data sensing head and an amplitude thereof exceeding a predetermined slice level and to generate, on the basis of the detected position of peak and the detected amplitude, a rectangular-wave read (reproduction) signal which represents the original or correct signal with high fidelity. As is well known, the peak position deviates from its original peak position due to a correlation between the peak position of interest and its adjacent peak positions (magnitude of the time interval). From this point of view, it is required to provide for an equalizing circuit which corrects a shift of the peak position (peak shift). Further, the amplitude of the reproduction signal varies due to a change of the reproduction frequency. Thus, it is required to provide for an equalizing circuit which corrects the amplitude variation. Moreover, the equalizing circuit directed to correcting a peak shift has an optimum value which is different from that for the equalizing circuit directed to correcting an the amplitude variation. For this reason, it is required to determine an optimum equalizing character, taking into consideration the difference between the optimum values of the two different equalizing circuits.\nReferring to FIG. 1, there is illustrated a conventional data reproducing circuit. The illustrated data reproducing circuit includes a data sensing head 10, which reads a recording medium (not shown) and generates a reproduction signal. A preamplifier 24 amplifies the reproduction signal from the head 10, and outputs an amplified reproduction signal. The amplified reproduction signal is pulled up to a power source voltage Vcc through a resistor 26, and is then supplied to a delay circuit 16-1 and an attenuator 18-1. An output from the delay circuit 16-1 is input to a non-inverting input terminal of a subtracter amplifier 20. An output from the attenuator 18-1 is input to an inverting input terminal of the subtracter amplifier 20, which subtracts the output supplied from the attenuator 18-1 from the output supplied from the delay amplifier 16-1, and outputs the result of this subtraction. The output from the subtracter amplifier 20 is pulled up to the power source voltage Vcc through a resistor 28, and is then input to a delay circuit 16-2 and an attenuator 18-2. An adder amplifier 22 adds an output from the delay circuit 16-2 and an output from the subtracter amplifier 18-2.\nThe illustrated data reproducing circuit has a first function of correcting a peak shift of the reproduction signal which occurs during reproduction and a second function of correcting an amplitude variation due to a change of the reproduction frequency. An optimum correction characteristic of the data reproducing circuit can be provided by adjusting the amount of attenuation in each of the attenuators 18-1 and 18-2.\nAs shown in FIG. 2, the reproduction signal supplied from the head 10 has negative edges 32 having amplitude components which have a polarity opposite to that of a main signal waveform 30 of the reproduction signal and which are located on both sides thereof. From this viewpoint, the attenuator 18-1 is adjusted so as to have a characteristic suitable for correcting the peak shift and amplitude variation of the main signal waveform 30, and the attenuator 18-2 is adjusted so as to have a characteristic suitable for correcting the peak shift and amplitude variation arising from the negative edges 32.\nIt is noted that generally the optimum values to be provided for the attenuators 18-1 and 18-2 for correcting the peak shift of the main signal waveform 30 are not equal to those to be provided for the attenuators 18-1 and 18-2 for correcting the amplitude variation thereof. For this reason, it is impossible to provide each of the attenuators 18-1 and 18-2 with an optimum value suitable for correcting both the peak shift and amplitude variation. For this reason, the data reproducing circuit shown in FIG. 1 cannot correct both the peak shift and the amplitude variation effectively and cannot reproduce the original signal with a high accuracy.\nThe above discussion holds true for another conventional data reproducing circuit as shown in FIG. 3, in which those parts which are the same as those shown in FIG. 1 are given the same reference numerals. The output from the preamplifier 24 is input to the delay circuit 16-1 and the attenuator 18-2. The output from the delay circuit 16-1 is input to the delay circuit 16-2 and the attenuator 18-1. The output from the delay circuit 16-2 is applied to a first non-inverting input terminal of an adder/subtracter amplifier 20'. The outputs from the attenuators 18-1 and 18-2 are applied to an inverting input terminal and and a second non-inverting input terminal of the adder/subtracter amplifier 20', respectively. The adder/subtracter amplifier 20' subtracts the output supplied from the delay circuit 16-2 from the output supplied from the attenuator 18-1, and adds the output from the attenuator 18-2 to the result of subtraction.\nThe arrangement shown in FIG. 3 has the same problem as the arrangement shown in FIG. 1. That is, the optimum values to be provided for the attenuators 18-1 and 18-2 for correcting the peak shift of the main signal waveform 30 are not equal to those to be provided for the attenuators 18-1 and 18-2 for correcting the amplitude variation. Thus, it is impossible to provide each of the attenuators 18-1 and 18-2 with an optimum value suitable for correcting both the peak shift and amplitude variation. For this reason, the data reproducing circuit shown in FIG. 3 cannot correct both the peak shift and the amplitude variation effectively and cannot reproduce the original signal with a high accuracy."}
{"text": "The invention relates to a drive unit for elevators consisting of a motor mounted in a casing, the motor driving a traction sheave via a drive shaft, there being a braking device to hold the traction sheave.\nFrom the German patent specification DD 44 278 an electric elevator motor has become known which has a traction sheave and brake drum which are fastened together. The traction sheave and brake drum are mounted outside both the motor casing and the bearing closest to the motor on the free end of a shaft in an overhung manner. The traction sheave and the brake drum can be constructed in one piece. The brake, which has two shoes, consists of an actuating device, brake arms, and brake shoes, the actuating device being fastened on the outside of the casing of the motor. To support the heavy loads occurring on the side closest to the traction sheave there is a bearing plate and a bearing for heavy loads.\nA disadvantage of the known device is that the overhung arrangement of the traction sheave and brake drum necessitate an expensive bearing plate and bearing for heavy loads. A further disadvantage is having the brake arms and brake shoes arranged at the sides, which increases the diameter of the drive unit.\nThe present invention concerns a drive unit with a narrow construction for use with an elevator.\nThe advantages resulting from the invention relate mainly to the fact that the drive unit is narrow in the direction at right angles to the drive shaft, and that the narrow construction makes it possible for the drive unit to be built into an elevator hoistway, for example in the space between the path of travel of the car and the wall of the hoistway. A further advantage lies in the common support of the traction sheave and drive shaft. There is no unsupported free end of a shaft, the traction sheave being supported directly on the bearing plate. As a result, bending forces on the drive shaft can be avoided."}
{"text": "1. Field of the Invention\nThe present invention relates to a collimator with an adjustable focal length, particularly in x-ray inspection systems.\n2. Description of the Background Art\nInspection methods with the use of x-rays are employed particularly in the detection of critical substances and objects in luggage or other freight. For this purpose, multi-stage systems are known whose first stage is based on the absorption of x-radiation. To detect certain critical substances such as, for example, explosives, a second stage is employed to which objects from the first stage are selectively supplied. Systems whose operating principle is based on diffraction phenomena are used as the second stage. In this case, the diffraction angle at which an incident x-ray is diffracted depends on the atomic lattice distance of the material to be analyzed and the energy and thereby the wavelength of the incident radiation. By analyzing the diffraction phenomena by means of x-ray detectors, conclusions can be reached about the lattice distance and thereby about the material. This type of two-stage system is disclosed, for example, in the German patent application 103305211.\nBecause x-ray inspection systems work with extremely low radiation intensities, highly sensitive detectors are employed. To avoid measurement inaccuracies, it must therefore be achieved that only the radiation produced by the testing device strikes the detector. In addition, care must be taken that radiation diffracted only at a single point is detected, because otherwise localization within the object to be examined is not possible. Spatial filtering is therefore necessary, which is performed by a so-called collimator.\nBecause it is technically very costly to generate monochromatic x-radiation, the highly limited x-ray used for testing, the so-called pencil beam, has an energy spectrum known, for example, from measurements. It follows from the Bragg equation that the incident radiation is diffracted at each point at an angle that depends on the radiation energy. Radiation with an energy spectrum is therefore diffracted within an angle range, and thereby the diffraction is rotation symmetric about the incident pencil beam. In an x-ray inspection, it is desirable to detect only radiation diffracted at a certain angle. This is also achieved with the use of a collimator. The passband of the collimator corresponds substantially to the generated surface of a cone whose tip coincides with the point whose diffraction properties are to be analyzed. To examine an area within an object, a plurality of points must be focused.\nGerman patent application 103305211 discloses a method for the examination of an object area in which the setup comprising a detector and collimator can be moved in the direction of the incident x-ray. The disadvantage of this method is that, on the one hand, a highly precise traveling unit is required and, on the other, the entire device must have an overall height more than twice the height of the object to be examined.\nA second possibility is the use of a collimator that has a plurality of parallel apertures of the same aperture angle and with which therefore a plurality of points on the rotation axis can be focused simultaneously. The use of a non-segmented detector, which is not position-sensing and therefore provides a common output signal for all focused points, however, results in the disadvantage that the evaluation and clear assignment of the detected radiation to a diffraction point are difficult. With use of a segmented detector, which, for example, is divided into separately evaluable circular rings, this disadvantage in fact does not arise, but this type of detector is laborious and costly."}
{"text": "The present invention relates to a socket for electrical parts for detachably accommodating and holding an electrical part such as a semiconductor device (called as xe2x80x9cIC packagexe2x80x9d hereinafter) and also relates to a method of assembling such socket for electrical parts.\nAs a conventional xe2x80x9csocket for electrical partsxe2x80x9d of this kind, there has been provided an IC socket for detachably accommodating and holding an IC package as xe2x80x9celectrical partxe2x80x9d.\nThe IC package includes, for example, a BGA (Ball Grid Array) type of IC package in which solder balls as a number of terminals are provided to the lower surface of the package body so as to protrude downward in a grid (lattice) arrangement having vertical (Y) and horizontal (X) rows.\nThe IC socket is provided with a socket body, a number of contact pins disposed to the socket body so as to contact the terminals of the IC package, respectively. Furthermore, a movable member for contacting or separating the contact pins to or from the solder balls of the IC package, through elastic deformation of the contact pins, is also arranged to be vertically movable.\nFurthermore, the socket body is provided with lever members to be rotatable (pivotal) about pivotal shafts so as to vertically move the movable member, and an operation member for rotating the lever members is also provided for the socket body to be vertically movable.\nWhen this operation member is vertically moved, the lever members are rotated (pivoted), and according to this rotating motion, the movable member is vertically moved, thus contacting or separating contact portions of the contact pins to or from the solder balls of the IC package.\nHowever, in such conventional IC package as mentioned above, a base end portion of the lever member is secured to the socket body to be rotatable through a rotation shaft and an E ring is also provided for preventing the lever member from coming off from the socket body, so that the number of elements or parts is increased and the assembling working is hence increased, thus being inconvenient and disadvantageous.\nFurthermore, in order to prevent the front end portion of the lever member from falling sideways, the front end portion of the lever member is coupled to the operation member to be rotatable through a support shaft, so that, in this meaning, the number of elements or parts is increased and the assembling working is hence increased, thus being inconvenient.\nAn object of the present invention is to substantially eliminate defects or drawbacks encountered in the prior art mentioned above and to provide a socket for electrical parts and a method of assembling the same, by which the number of elements or parts to be utilized can be effectively reduced as well as reduction of the assembling workings.\nThis and other objects can be achieved according to the present invention by providing, in one aspect, a socket for electrical parts in which a number of contact pins contacting a number of terminals of an electrical part is arranged to a socket body in which the electrical parts is accommodated, an operation member is provided for the socket body to be vertically movable in an installed state, and when the operation member is vertically moved, a movable member is moved through lever members to thereby displace the contact pins to be electrically contacted to or separated from the terminals of the electrical part,\nwherein the socket body is formed with boss portions to which base end portions of the lever members are secured to be rotatable and the socket body is formed with guide wall sections so as to guide side portions of the base end portions of the lever members to thereby prevent the lever members from coming off from the boss portions.\nIn a preferred embodiment of the above aspect, the operation member is formed with guide grooves into which front end portions of the lever members are inserted so as to prevent the lever members from falling sideways.\nThe socket body is formed of synthetic resin, the boss portions are integrally formed with the socket body and the lever members are formed of metal material. The base end portions of the lever members are formed with holes into which the boss portions are fitted. The terminals of the electrical part are solder balls. The socket is an IC socket and the electrical part is an IC package.\nIn a modified aspect, there is provided a socket for electrical parts comprising:\na socket body to which an electrical part having a number of terminals is accommodated;\na number of contact pins arranged to the socket body so as to accord with the arrangement of the terminals of the electrical part;\nan operation member disposed to the socket body to be vertically movable in an installed state;\nlever members provided for the socket body to be rotatable in association with the vertical movement of the operation member; and\na movable member disposed to the socket body to be movable in accordance with the movement of the operation member through the lever members to thereby displace the contact pins so as to contact or separate the contact pins to or from the terminals of the electrical part,\nthe socket body being formed with boss portions, the lever members being formed with base end portions which are fitted to the boss portions to be rotatable, and the socket body being formed with guide wall sections so as to guide side portions of the base end portions of the lever members to thereby prevent the lever members from coming off from the boss portions.\nIn another aspect, there is also provided a method of assembling a socket for electrical parts in which a number of contact pins contacting a number of terminals of an electrical part is arranged to a socket body which accommodates the electrical part and which is provided with an operation member and a movable member which is operatively connected to the operation member through lever members,\nwherein the lever members are fitted to the boss portions of the socket body to be rotatable, when the lever members are rotated while being guided along guide wall sections of the socket body and the operation member is disposed to the socket body to be vertically movable.\nThe operation member is formed with guide grooves into which front end portions of the lever members are inserted so as to prevent the lever members from falling down sideways.\nAccording to the present invention of the characters mentioned above, the base end portions of the lever members are fitted to the boss portions of the socket body to be rotatable and the guide wall sections are formed to the socket body so as to guide the side portions of the base end portions of the lever members to prevent the lever member from coming-off from the boss portions. Accordingly, it is not necessary to specifically locate any member for preventing the coming off of the lever members, thus reducing the number of parts or elements to be assembled and improving the assembling working of the lever members and others. Furthermore, the formation of the guide grooves eliminates the location of any member such as a support shaft for supporting the lever members.\nMoreover, the socket body is formed of synthetic resin and the lever members are formed of metal material. Accordingly, the boss portions can be easily molded integrally with the socket body and the lever members can ensure the good rigidity thereof.\nThe nature and further characteristic features of the present invention will be made more clear from the following descriptions made with reference to the accompanying drawings."}
{"text": "1. Field of the invention\nThe present invention concerns processes that allow measuring the induced magnetization in a nautical vessel in order to be able to thereafter carry out a process for the magnetic immunization of said vessels. It also concerns devices that allow this process to be carried out.\n2. Summary of the prior art\nAt the present time, most nautical vessels, boats and submarines are built from steel or iron and thus have a high magnetization. This magnetization comprises, on the one hand, a permanent proportion that is specific to vessels and, on the other hand, a variable proportion induced by the magnetic field of the earth and which thus depends upon the position of the vessel with respect to the magnetic-field of the earth.\nThis magnetization of the vessel is superimposed upon a magnetic field of the earth and thus generates a disturbance called the \"magnetic signature\" of the vessel. This disturbance allows for the location of the vessel to be and determined thereby making it possible to guide or fire missiles intended to destroy it. It is therefore extremely important to minimize, or even suppress, this disturbance in order to avoid the vessel detected by a magnetic method.\nThis operation, called \"magnetic immunization\", is performed in a manner known per se by creating in the volume of the vessel a magnetic field that is opposed to that of the vessel in order to cancel the magnetic signature.\nTo do this,the vessel is fitted with an assembly of circuits known as \"immunization loops\" through which pass electric currents. The sizes and the disposition of the loops, and the intensity of the currents that pass through them, are determined so as to fully minimize the magnetic signature of the vessel, whatever its orientation in the magnetic field of the earth, i.e. whatever its head and its list due to rolling and to pitching. These immunization loops are distributed throughout three assemblies corresponding to the axes of roll, pitch and rocking, and conventionally called \"L, M, A\".\nTo determine the intensities of the currents that will pass through the immunization loops, it is necessary to measure the magnetic signature of the vessel and for this purpose a magnetic measuring station is used.\nA magnetic measuring station advantageously comprises two linear networks of magnetic sensors, called bases, placed on the sea-bed and each aligned according to a cardinal direction, N/S for one, and E/W for the other. These sensors are connected to a land-based processing unit which receives the signals produced by the passage of the vessel to be immunized above the bases according to trajectories which are preferably themselves cardinal. These signals are processed in the station to determine the intensity and the polarity of the currents in the immunization loops in such a way as to obtain satisfactory immunization, whatever the orientation of the vessel in the magnetic field of the earth.\nSince the permanent magnetization is steady with respect to the vessel in direction and amplitude and only evolves very slowly as a function of time, it is possible to determine a continuous component of the current in each loop, the value of which will be set and possibly reset during a subsequent immunization operation.\nOn the other hand, the induced magnetization is variable and it is necessary to superimpose upon this continuous component a variable component determined according to the heading and the orientation of the vessel which are known by measuring means, such as gyroscopic or optical means, for example.\nTo determine respectively these continuous components and variable components, it is necessary, when measuring the magnetic signature to separate the influences of the permanent magnetization from those of the induced magnetization in the total magnetization, and this separation must be made in the three directions corresponding to the three axes of the vessel.\nTo do this, the vessel is usually required to travel along the same route twice above the bases according to opposite headings. The vessel thus passes in a N/S direction above the E/W base and then S/N above this same base. It thereafter passes in E/W direction above the N/S base and returns in a W/E direction above said same base. The orientation of the bases and the routes is not compulsory but it facilitates the interpretions measurements and calculations.\nDuring these opposite passages, the permanent magnetization, which is associated to the vessel, turns with the vessel while the induced magnetization does not turn. In order to determine this induced magnetization, it is thus sufficient to subtract the two measurements corresponding to two passages in opposite directions. Once the induced magnetization and the total magnetization are known, the permanent magnetization is directly obtained.\nThe fact of having to complete two successive passages according to opposite heads is in itself a drawback, especially since it lengthens the duration of operations. Furthermore, it is necessary to set back the measurements, for example by interpolation, so as to perform subtraction on homologue points for both passages. Indeed, an inevitable variation exists between the routes in one direction and in the other. Since, furthermore, the trajectory measurements themselves are inaccurate, it can be seen that sources of errors accumulate.\nSimilarly, this method is limited to longitudinal and transversal magnetization, since it is obviously not possible to set the vessel vertical in order to separate the induced magnetization from the permanent magnetization according to the vertical direction. In this latter case, the two magnetizations are separated according to empirical methods based on the experience of the operators and measurements on models that are roughly representative of the vessel.\nIn order to overcome these drawbacks, the invention proposes to create in the vessel a supplementary known magnetic field that allows to generate a supplementary induced magnetization, the resulting magnetic field of which can be easily separated from that due to the magnetization of the magnetic field of the earth. This supplementary magnetic field is calculated so as to be geometrically as close as possible to that of the magnetic field of the earth in the vessel. It is obtained by means of alternating currents injected into the immunization loops. The magnetic field produced by the loops along with respect to the sensors can be calculated, thereby allowing to obtain by subtraction from the alternating magnetic field measured by the sensors of the immunization base the magnetic field due to the induced field in the vessel."}
{"text": "In many applications it is important to find objects that are carried by people and which can be concealed on their bodies. Detection systems are known which require an individual to pass through a fixed detector passageway (portal) equipped with magnetic sensors. When metallic objects of sufficient mass pass through the passageway, a warning signal is activated because a change in magnetic flux is detected. This type of system either detects or does not detect a metal object and sometimes makes no determination relative to the amount of metal present. Keys, jewelry, watches, and metal-framed eyeglasses may all trigger such a system.\nWhile such magnetic detectors are capable of detecting metal objects passing through the detector passageway, such systems cannot determine whether the detected metal object is a threat object (e.g., a knife or a gun) or an innocuous (non-threat) object (e.g., keys, coins, jewelry, belt buckles). Moreover, these detection systems do not pinpoint the location of the metal objects on the individual's body. Likewise, these systems are useless in detection of modern threats posed by plastic and ceramic items and plastic and liquid explosives.\nMoreover, there are situations when a person under surveillance need not be aware that he or she is being monitored. In such a case, the portal systems described above are not appropriate.\nAnother type of detection system also employs imaging techniques that acquire images of the detection space and then display the image to an operator. Moreover, imaging detection systems can use image recognition methods to convert the image into an indication (such as an audible or visual alarm). In order to recognize a specific threat object, the system has to compare the object with an electronic catalog of images of uniquely-shaped threat objects. In this case, unique orientations of the objects are also important, because an object may have a significantly different appearance if viewed from the sides, the top, etc. The observed uniqueness of a threat object also essentially depends on the image resolution of the system.\nA good example of a concealed object imaging system that exhibits high image clarity is the cabinet x-ray system used at airports to screen carry-on luggage. Although very effective for certain security tasks, X-ray imaging can pose a serious health risk to living organisms due to X-ray exposure, and is therefore unacceptable to the public. On the other hand, RF radiation in the microwave and millimeter-wave (e.g., 5 GHz to 1 THz) range offers a possible solution for concealed weapon detection and imaging, because the RF radiation can easily penetrate clothing and also represents no known health threat to humans at moderate power levels (see, for example, U.S. Pat. No. 6,791,487 to Singh; U.S. Pat. No. 6,876,322 to Keller; U.S. Pat. No. 6,992,616 to Grudkowski et al; and U.S. Pat. No. 6,965,340 to Maharav at al.).\nSpecifically, U.S. Pat. No. 6,791,487 describes imaging methods and systems for concealed weapon detection. In an active mode, a target can be illuminated by a wide-band RF source. A mechanically scanned antenna, together with a highly sensitive wide-band receiver can then collect and process the signals reflected from the target. In a passive mode, the wide-band receiver detects back-body radiation emanating from the target and possesses sufficient resolution to separate different objects. The received signals can then be processed via a computer and displayed on a display unit thereof for further analysis by security personnel.\nA wideband millimeter-wave imaging system is described in the article titled “Concealed explosive detection on personnel using a wideband holographic millimeter-wave imaging system,” by Sheen et al., Proceedings of SPIE—The International Society for Optical Engineering, V. 2755, 1996, PP. 503-513. To form an image, Sheen et al. use a linear array of 128 antennas that can electronically scan over a horizontal aperture of 0.75 meters, while the linear array is mechanically swept over a vertical aperture of 2 meters. At each point over this 2-D aperture, coherent wideband data reflected from the target is gathered using wide-beamwidth antennas. The data is recorded coherently, and reconstructed (focused) using an image reconstruction algorithm that works in the near-field of both the target and the scanned aperture and preserves the diffraction limited resolution of less than one-wavelength. The wide frequency bandwidth is used to provide depth resolution, which allows the image to be fully focused over a wide range of depths, resulting in a full 3-D image.\nU.S. Pat. No. 6,965,340 describes a security inspection system including a portal through which a human subject is capable of walking and a scanning panel including an array of antennas that are programmable with a respective phase delay to direct a beam of microwave illumination toward a target on the human subject. The antennas are further capable of receiving reflected microwave illumination reflected from the target. A processor is operable to measure an intensity of the reflected microwave illumination to determine a value of a pixel within an image of the human subject. Multiple beams can be directed towards the human subject to obtain corresponding pixel values for use by the processor in constructing the image.\nU.S. Pat. No. 6,992,616 describes active imaging systems including an antenna apparatus configured to transmit toward and receive from a subject in a subject position, electromagnetic radiation. The antenna apparatus may move in a partial or continuous loop around the subject, toward or away from the subject, or in an opposite direction to an associated antenna apparatus. Antenna units in the antenna apparatus may be oriented at different angular positions along an array. Antenna arrays may also be formed of a plurality of array segments, and a group of arrays may be combined to form an antenna apparatus."}
{"text": "Wheeled carriages for supporting a patient in a substantially horizontal position are well-known in the art and a representative example of an early version of such a device is illustrated in Dr. Homer H. Stryker's U.S. Pat. No. 3,304,116, reference to which is incorporated herein. Dr. Stryker's innovative wheeled carriage included a fifth wheel which is raisable and lowerable by the attendant by directly manually manipulating the wheel support frame oriented beneath the patient supporting portion of the wheeled carriage. The orientation of the fifth wheel was sometimes awkward to reach and, therefore, made the operation of the raising and lowering feature of the fifth wheel difficult to attain.\nOther structure was added to the wheeled carriage to facilitate an activation of the brakes for the wheels on the wheeled carriage from positions adjacent the head end and/or the foot end of the wheeled carriage. However, if the wheeled carriage were to be placed into a position where the head end and the foot end of the wheeled carriage were inaccessible to the attendant, operation of the brake became difficult without first moving the wheeled carriage to a position wherein at least one of the head and/or foot end of the wheeled carriage would be accessible for operation of the brake. If a fifth wheel is present and is deployed to its floor engaging position, situations where this might be considered a problem would be where an overbed table was to be placed in association with the wheeled carriage and the fifth wheel was blocking entry of the wheeled carriage of the overbed table beneath the wheeled carriage because of the presence of the lowered fifth wheel. Thus, it became a desire to provide an easily accessible fifth wheel and brake activation device oriented at least within the lateral side region of the wheeled carriage as well as within the head and foot regions of the wheeled carriage.\nAs wheeled carriages for supporting a patient further developed from Dr. Stryker's earlier patent, the mechanism for raising the patient support relative to the wheeled base generally included a pair of horizontally spaced hydraulic jacks which were simultaneously pumped with hydraulic fluid by operation of a single foot activated pedal. Once the hydraulic jacks had raised the patient support to the desired elevation, either the head end of the patient support, the foot end of the patient support or both ends of the patient support could be selectively lowered by activation of one or two foot activated pedals. For example, one foot activated pedal, when depressed, would activate a hydraulic fluid release valve for allowing hydraulic fluid to exit the hydraulic jack at one end of the bed so that that end of the bed would be lowered. The second foot pedal would accomplish the same task. When it was desired to lower both the head end and the foot end of the patient support at the same time, it was necessary for both foot pedals to be depressed at the same time. Attendants have found this difficult to achieve. Accordingly, it became a desire to provide for an easy to use mechanism for effecting the simultaneous lowering of the head end and foot end hydraulic jacks.\nAccordingly, it is an object of this invention to provide a wheeled carriage for supporting a patient in a substantially horizontal position having a wheel braking and unbraking mechanism and/or an auxiliary wheel and support structure therefor mounted on a wheeled base, one and/or the other being actuatable by a manually manipulatable control element at at least one of the pair of lateral side regions or at least one of the head or foot ends of the wheeled carriage so that an attendant can operate the manually manipulatable control element to effect a movement of the auxiliary wheel solely from the head or foot end and solely from within the lateral side region.\nIt is a further object of this invention to provide brakes for the wheels of the wheeled carriage and a control mechanism for activating the brakes while the auxiliary wheel is in a position spaced from the floor surface and deactivating the brakes while the auxiliary wheel is in a floor engaging position, all utilizing the aforesaid same control mechanism.\nIt is a further object of this invention to provide a wheeled carriage, as aforesaid, wherein plural control elements are provided around the perimeter of the wheeled carriage to facilitate an attendant operating a selected one of the manually manipulatable control elements to effect a movement of the auxiliary wheel from its raised or lowered position and/or activation of a brake mechanism for the wheeled carriage solely from within a selected one of the head, foot and two lateral side regions of the wheeled carriage.\nIt is a further object of the invention to provide a control mechanism for actuating the raising and lowering feature of the fifth wheel and/or activation of a brake mechanism for the wheeled carriage by utilizing a rotational movement of the activating devices to facilitate compact construction of a rotary transmission device to interconnect the multiple locations for activating the raising and lowering of the fifth wheel feature and/or activation of a brake mechanism for the wheeled carriage.\nIt is a further object of the invention to provide a wheeled carriage, as aforesaid, wherein the manually manipulatable control element at each of the multiple locations around the perimeter of the wheeled carriage are identical to one another thereby standardizing the appearance of the control element to the attendant thereby minimizing confusion as to which of the many manually manipulatable elements on a wheeled carriage for supporting a patient in a substantially horizontal position is to be activated.\nIt is a further object of the invention to provide a wheeled carriage, as aforesaid, wherein the fifth wheel activating structure is durable and requires little or no maintenance over the lifetime of the wheeled carriage.\nIt is a further object of the invention to provide a wheeled carriage, as aforesaid, wherein the control element for activating the brakes and/or the auxiliary fifth wheel is a unitary pedal construction.\nIt is a further object of the invention to provide a wheeled carriage, as aforesaid, wherein hydraulic jacks are utilized to raise and lower the patient support relative to the wheeled base and wherein a unitary pedal construction is utilized to effect an independent lowering of the head end and the foot end of the patient support as well as a simultaneous lowering of both the head end and the foot end of the patient support."}
{"text": "Ever since the first experiments in variolation in 1721, and Jenner's vaccination methods in 1796, methods and compositions for disease prevention utilizing immunization have been extensively investigated. Many methods rely upon the use of active immunization, in which an antigen (or mixtures of antigens), such as a modified infectious agent or toxin is administered, resulting in active immunity. This active immunity is characterized by the production of antibodies directed against the administered antigen(s), and in some cases, induction of cellular responses mediated by lymphocytes and macrophages.\nTraditionally, vaccines used for active immunization have consisted of live attenuated bacteria (e.g., Bacillus Calmette-Guerin) or viruses (e.g., measles virus), killed microorganisms (e.g., Vibrio cholerae), inactivated bacterial products (e.g., tetanus toxoid), or specific single components of bacteria (e.g., Haemophilus influenzae polysaccharide). Although active immunization with live organisms is generally superior to immunization with killed vaccines in producing long-lived immune responses, care must be taken to properly store and administer these vaccines, as serious failures of measles and smallpox immunizations have resulted from improper refrigeration of the vaccine preparations. In addition, pregnant women and individuals with compromised immune systems should, in general, not receive live vaccines, as the organisms may cause serious disease upon vaccination. For example, live vaccines have caused serious and fatal disease in patients receiving corticosteroids, alkylating drugs, radiation, other immunosuppressive treatments, as well as individuals with known or suspected congenital or acquired defects in cell-mediated immunity (e.g., severe combined immunodeficiency disease, leukemia, lymphoma, Hodgkin's disease, and acquired immunodeficiency syndrome [AIDS]). Live vaccines may even cause mild, or rarely, severe disease in immunocompetent hosts. In addition, live vaccines may also contain undesirable components. For example, epidemic hepatitis has resulted from the use of vaccinia and yellow fever vaccines containing human serum.\nPassive immunization using preformed immunoreactive serum or cells is sometimes utilized, especially when active immunization is not available or not advisable. In particular, passive immunization finds use in individuals who cannot produce antibodies or other immune system deficiencies, as well as in individuals who are at risk of developing disease before active immunization would be successful in stimulating a sufficient antibody response. Passive immunization is also used in conjunction with vaccine administration in the management of certain diseases (e.g., rabies vaccination and prophylaxis following an animal bite), management of individuals who have been exposed to certain toxins or venoms, and as an immunosuppressant. However, passive immunization does not produce long-term immunity and is sometimes associated with severe reactions due to the presence of foreign proteins in the vaccine preparation (e.g., anaphylaxis resulting from a reaction against human or horse [or other non-human animal] proteins present in the vaccine preparation).\nMore recently, vaccines comprising recombinant DNA or RNA segments have been developed. However, use of these recombinant vaccines has resulted in problems associated with the expression of the desired antigen(s) in another organism (e.g., an E. coli or yeast host). For example, in addition to the desired antigen, other components, such as other antigens (e.g., protein and other components) from the expression host, preservatives, etc may be present in the preparation. In addition, adjuvants are sometimes required in order to provide efficacious vaccination with these vaccines. However as with passive immunization, undesirable reactions sometimes occur in vaccinated individuals due to the presence of these undesirable components.\nVarious adenovirus-based gene delivery systems have likewise been investigated for vaccine use. Human adenoviruses are double-stranded DNA viruses which enter cells by receptor-mediated endocytosis. These viruses have been viewed as being particularly well suited for gene transfer because they are easy to grow and manipulate and they exhibit a broad host range in vivo and in vitro. Adenovirus is easily produced at high titers and is stable so that it can be purified and stored. Even in the replication-competent form, adenoviruses generally cause only low level morbidity and are not associated with human malignancies. Various references provide reviews of adenovirus-based gene delivery systems (See, e.g., Haj-Ahmad and Graham, J. Virol., 57:267-274 [1986]; Bett et al., J. Virol., 67:5911-5921 [1993]; Mittereder et al., Human Gene Ther., 5:717-729 [1994]; Seth et al., J. Virol., 68:933-940 [1994]; Barr et al., Gene Ther., 1:51-58 [1994]; Berkner, BioTechn., 6:616-629 [1988]; and Rich et al., Human Gene Ther., 4:461-476 [1993]). However, despite these advantages, adenovirus vector systems still have several drawbacks which limit their effectiveness in gene delivery, such as cytotoxicity. Adenovirus vectors also express viral proteins that may elicit a strong non-specific immune response in the host. This non-specific immune reaction may increase toxicity or preclude subsequent treatments because of humoral and/or T cell responses against the adenoviral particles. Thus, problems remain even with the newer technologies for vaccine administration.\nAs briefly mentioned above, the major focus in the past has been on the development of antibody responses to vaccination. However, cell-mediated responses are of great importance in some situations. Indeed, cell-mediated immunity is of greater importance than the antibody-mediated response in the response to intracellular parasites (e.g., viruses and obligately intracellular bacteria). T-cells (T lymphocytes) play the primary roles in cell-mediated immunity, although there is communication via cytokines and other signalling compounds between these cells as the antibody-producing B-cells.\nCytotoxic T-lymphocytes (CTLs) play an important role in immune responses directed against intracellular pathogens such as viruses and tumor-specific antigens produced by cancerous cells. In particular, CTLs mediate cytotoxicity of virally infected cells by recognizing viral determinants in conjunction with Class I MHC molecules displayed by the infected cells. Cytoplasmic expression of proteins is a prerequisite for Class I MHC processing and presentation of antigenic peptides to CTLs. However, conventional immunization techniques, such as those using killed or attenuated viruses, often fail to elicit an appropriate CTL response which is effective against an intracellular infection. Thus, there remains a need for the development of vaccines that stimulate appropriate responses (i.e., cell-mediated as well as antibody-mediated immune responses), in order to prevent disease. Indeed, despite advances in vaccine technology, there remains a need for vaccines that are efficacious, yet avoid the problems associated with current vaccine preparations."}
{"text": "Due to advances in polymer science, the number of applications for polymeric tubing have drastically increased. These advances have resulted in polymeric tubing having outstanding thermal, mechanical, and insulative properties. Moreover, due to the wide variety of base resins, additives, and processing techniques currently available, tubing can be designed having specific physical properties, e.g. flexibility, flame resistant, chemical resistant, etc.\nTypically, the base resins used for such polymeric tubing include polyolefins, polyvinylchloride (PVC), fluoropolmers, elastomers, and blends thereof.\nDepending upon the materials and processing techniques used, such tubing may be \"heat-shrinkable\". For example, polyolefin materials are commonly extruded, irradiation crosslinked, and then expanded to form heat shrinkable tubing, as is common in the art and as described in: R. Kraus and D. Ryan, \"Advances in Heat-Shrink Technology,\" IEEE Electrical Insulation Magazine (1988) Vol.4, 31-34; and J. W. Hoffman, \"Insulation Enhancement with Heat-Shrinkable Components,\" IEEE Electrical Insulation Magazine (1991) Vol.7, 33-38. Upon subsequent application of heat, such tubing shrinks to approximately its originally extruded size and shape.\nAlthough heat shrinkable polyolefin tubing exhibits many desirable mechanical, thermal, and insulative properties, in many applications including those having repeated exposure to high temperature and/or pressure, such polyolefin materials fail to maintain an adhesive bond to a substrate and/or fail to maintain heat sealability (ability to remain bonded to itself). Consequently, in these applications, an inner sealant or adhesive mastic liner is commonly used to help maintain a fluid tight seal. For example, U.S. Pat. No. 3,297,819 to Wetmore; R. Kraus and D. Ryan, \"Advances in Heat-Shrink Technology,\" IEEE Electrical Insulation Magazine (1988) Vol.4, 31-34; and J. W. Hoffman, \"Insulation Enhancement with Heat-Shrinkable Components,\" IEEE Electrical Insulation Magazine (1991) Vol.7, 33-38 all describe a heat shrinkable tubing including an inner thermoplastic or mastic liner which maintains an adhesive bond with an inner substrate, and an outer heat shrinkable liner which is typically crosslinked. Such tubing may be co-extruded thereby providing a tubing having the inner non-crosslinked liner and the outer crosslinked liner.\nMany problems are associated with such tubing, particularly when used in environments having cyclic exposure to high temperatures and pressures. For example, when such tubing is used in the medical or food industries where cyclic autoclave sterilization is common, the inner liner of the tubing melts and flows in response to the sterilization temperatures and pressures. Moreover, at sterilization temperatures, the outer liner tends to undergo additional shrink thereby causing the melted inner liner to ooze or flow from the tubing. The melt and flow of the inner liner, although acceptable in some applications, is not acceptable in most medical applications. For example, one particular medical applications includes the use of the heat shrinkable tubing in connection with electro-surgical laparoscopic instruments. Such instruments typically include a cylindrical electrically conducting member having one of many possible surgical attachments secured to one end for performing a variety of surgical procedures, i.e. providing suction, irrigation, coagulating vessels, etc. The opposite end of the conducting member is securable to a hand-held control module which allows a surgeon to control the surgical attachment.\nThe conducting member includes a sheath or tubing disposed circumferencally about its length for electrically insulating the conducting member. The tubing is preferably transparent and of the heat shrink variety so that it may be easily applied about the conducting member. Moreover, the tubing must maintain an adhesive seal with the conducting member and maintain heat sealability with itself after cyclic autoclave sterilization in order to prevent the ingress of moisture between the conducting member and the tubing. The tubing must maintain its electrical insulating properties along with a sufficient hot modulus. Furthermore, the tubing must maintain all of these aforementioned properties after cyclic gamma, electron beam, or ethylene oxide gas sterilization procedures. Finally, the tubing must maintain its thermal stability at sterilization temperatures and not flow in response to sterilization heat and pressure.\nPrior art heat shrink polyolefin tubing, including those made from KYPRAN.TM. (registered trademark of Penwalt Co. for its vinylidene fluoride resin) and crosslinked polyolefins such as those disclosed in U.S. Pat. No. 3,592,881 to Ostapchenko, U.S. Pat No. 3,990,479 to Stine et al., and U.S. Pat. No. 3,852,177 to Atchison et al., may provide the necessary hot modulus strength, but lose their adhesive sealability to a substrate and lost sealability after exposure to the high temperatures and pressures associated with autoclave sterilization. Moisture leaks into the interface between the conducting member and the tubing, thus causing potential electrical and sterilization problems. Such prior art tubing also permanently discolor after cyclic gamma sterilization. Furthermore, the prior art tubing lose their electrical insulating properties after cyclic sterilization. Although co-extruded tubing such as that shown in U.S. Pat. No. 3,297,819 to Wetmore, may maintain a seal after repetitive autoclave sterilizations, the temperature and pressure associated with sterilization cause the inner liner of adhesive material to flow out of the tubing rendering the instrument unusable and susceptible to leaks.\nThus, a tubing is needed which is thermally stable and will not flow at sterilizations temperatures while simultaneously providing sufficient hot modulus, electrical insulative properties, permanent transparency, heat sealability, and adhesiveness after cyclic autoclave, gamma, and electron beam sterilization."}
{"text": "An electrical power system operates under a steady-state condition when there exists a balance between generated and consumed active power for the system. Power system disturbances may cause oscillations in machine rotor angles that can result in conditions like a power swing, when internal voltages of system generators slip relative to each other. Power system faults, line switching, generator disconnection, or the loss or sudden application of large amounts of load are examples of system disturbances that may cause a power swing event to occur in a power system. Depending on the severity of the disturbance and power system control actions, the system may return to a stable state or experience a large separation of load angle and eventually lose synchronism. Large power swings, stable or unstable, may cause unwanted relay operations at different locations in the system, which can aggravate the system disturbance and can result in major power outages or blackouts.\nFurther, asynchronous operation of interconnected generators in the power system as an effect of unstable power swing may initiate uncontrolled tripping of circuit breakers resulting in equipment damage and posing a safety concern for utility operators. Therefore, the asynchronous system areas may need to be separated from each other quickly and dynamically in order to avoid extensive equipment damage and shutdown of major portions of the system. In order to contain these risks, it is required as per international standards to have an optimal generator protection device, such as a generator relay, in place to isolate generators from the rest of the system within a half-slip cycle. The need to meet the international standards challenges protection engineers to ensure selective and reliable relay operation.\nIn a conventional relaying approach, a variation in system impedance determined at generator terminals is analyzed for detecting power swing. Various impedance-based protection approaches including power swing block (PSB) and out-of step trip (OST) are currently being used. However, these protection approaches may need an extensive power system stability study to arrive at an optimal setting for selective and reliable relay operation. Protection engineers typically use preliminary settings that are not adapted to accommodate variation in system configurations or operational dynamics, for example, changes in transmission and distribution layout during implementation phase or dynamically during operational phase. Extensive study and non-dynamic preliminary settings may result in the protection device being unable to selectively, reliably and dependably detect power swings and isolate generators during such events.\nOther known relaying approaches estimate swing center voltage (SCV) for detecting power swings. Such approaches use approximate estimation that does not take into consideration real time power system dynamics. In some relaying approaches, a high-speed communication network such as fiber optic or global positioning system (GPS) communication is used to obtain data at a source end from one or more generators at receiving end(s), which is at a remote location from the source end, for SCV estimation. However, such approaches have economic challenges due to the cost associated with implementing and maintaining a high-speed communication network. Some approaches for SCV directly measure the rotor angle between the generator's internal voltage and terminal voltage for detecting power swing. In the absence of direct measurements, it is difficult to determine the power swing condition.\nIn one known SCV approach, the relationship between the SCV and a swing angle (θ) of a two-source equivalent system may be determined as per the below equation:\n  SVC  =            ±      E        ×          cos      ⁡              (                  θ          2                )            where, E is an internal voltage of a source-end generator\nIn such approaches, the power swing may be detected by calculating a rate of change of the SCV. The time derivative of the SCV is given by below equation:\n            d      ⁡              (        SCV        )              dt    =            E      2        ×          sin      ⁡              (                  θ          2                )              ⁢                  d        ⁢                                  ⁢        θ            dt      \nIn this equation, for sin (θ/2) to be close to one, θ should be around 180 degrees (for example, between 90 and 180 degrees). Therefore, the above equation can be used for detecting power swing when the value of θ is around 180 degrees. However, for values of θ between 0 and 90 degrees, the above equation will result in sin (θ/2) to be close to zero. In other words, this approach is not suitable for a smaller range of values of θ (for example, between 0 and 90 degrees)."}
{"text": "Aromatic compounds can be alkylated with olefinically unsaturated compounds employing a catalyst consisting essentially of aluminum halide and molecular iodine. Additionally, molecular iodine or iodine-containing compounds are useful in other catalytic conversions such as for example, Grignard reactions, carbonylation reactions, and the like. A problem encountered with such reactions is the removal of the iodine components from the reaction mixture so that color formation in the product is avoided. Methods disclosed in the prior art, such as for example, extraction with thiosulfate or bisulfite; treatment with other oxidants; distillation in the presence of acid or base, however, are not always successful in completely removing color forming impurities from the reaction mixture which contains the desired product."}
{"text": "Most of the detergents in use today are derived from precursor petrochemicals. The currently predominant precursors is linear alkyl benzene (LAB), which is commonly produced by the alkylation of benzene with a long chain linear olefin. The subject invention is directed to the production of monomethyl acyclic olefins and paraffins, which may be recovered as a product in their own right, or used in the production of various petrochemicals as through alkylation or oxygenation. The following description of the invention will mainly address the recovery and use of the monomethyl hydrocarbons in the production of detergent precursor petrochemicals, and in particular the production of alkylbenzene derived detergents.\nSeveral quality characteristics of alkylbenzenesulfonate (ABS) detergents are set by the chemical structure of the alkyl side chain. For instance, linear alkyl groups have the advantage of increased biodegradability. Other characteristics of the detergent such as its effectiveness in hard water and its foaming tendency are also influenced by the structure of the side chain and its constituents. It has recently been determined that highly desirable detergent precursors can be formed from olefins which contain a single methyl side chain on the main alkane chain. This is a departure from the previous preference for straight chain alkanes. The subject invention is specifically directed to the production of monomethyl hydrocarbons for use in the subsequent production of these detergent precursors or ingredients."}
{"text": "The invention relates to a driver assistance device having a plurality of ultrasound sensors. Moreover, the invention also relates to a vehicle having such a driver assistance device as well as a method for operating a driver assistance device.\nA plurality of different driver assistance systems or driver assistance devices for motor vehicles is known. A number of such driver assistance devices operate on the basis of ultrasound sensors that are designed for detecting the surroundings of a vehicle. Specific actions are carried out depending on this information, which is detected by means of the ultrasound sensors, e.g. information is output to the driver of the vehicle and/or at least one semi-autonomous process is carried out for driving the vehicle by the driver assistance device.\nWith current parking assistance systems or parking assistance devices the distance between an object and the vehicle is determined using the transition time of previously transmitted ultrasound signals. An ultrasound sensor transmits a defined number of ultrasound signals, which in turn are received by one or a plurality of the ultrasound sensors. These are always the same ultrasound signals. It is therefore not possible for adjacent ultrasound sensors to transmit simultaneously, because the received ultrasound signals cannot otherwise be unambiguously assigned.\nMoreover, a method of operating an ultrasound sensor array is known from DE 101 06 142 A1. Said array comprises at least two transmitting units distributed along the circumference of the motor vehicle for transmitting ultrasound pulses, and at least one receiving unit for receiving ultrasound pulses reflected at an object in the monitoring area. A plurality of transmitting units can be operated in parallel and at the same time transmit ultrasound pulses that are encoded relative to each other. For encoding, carrier signals of the ultrasound pulses are frequency modulated differently for the individual simultaneously operated transmitting units at least during the pulse period and the received encoded ultrasound pulses are assigned to individual simultaneously operated transmitting units by the at least one receiving unit using the encoding.\nRegarding the encoding of ultrasound signals, a plurality of options and procedures are mentioned in said DE 101 06 142 A1. Besides frequency encoding, in which the frequency of the carrier signal of a first transmitting unit is upward modulated and the frequency of the carrier signal of a second transmitting unit is downward modulated, amplitude modulation can also be carried out over one or a plurality of pulses.\nSuch ultrasound sensors, which transmit encoded ultrasound signals, are however relatively expensive.\nMoreover, it is not necessary for each function of a driver assistance system to transmit such encoded ultrasound signals, but conventional unencoded signals are also sufficient.\nThe object of the present invention is to provide a driver assistance device with whose ultrasound sensors high functionality in respect of safe and reliable detection of the surroundings is guaranteed and which can be implemented inexpensively.\nThis object is achieved by a driver assistance device, a vehicle having such a driver assistance device and a method for operating such a driver assistance device according to the independent claims."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a gas monitoring system, and, in particular, to an integrated sample cell and filter for use in a sidestream gas monitoring system, such filter for separating undesired liquid components from respiratory gases to be monitored in the sample cell.\n2. Description of the Related Art\nDuring medical treatment, it is often desirable to monitor and analyze a patient's exhalations to determine the gaseous composition of the exhalate. For instance, monitoring the carbon dioxide (CO2) content of a patient's exhalations is often desirable. Typically, the carbon dioxide (or other gaseous) content of a patient's exhalation is monitored by transferring a portion, or sample, of the patient's expired gases to a suitable sensing mechanism and monitoring system.\nMonitoring of exhaled gases may be accomplished utilizing either mainstream or sidestream monitoring systems. In a mainstream monitoring system, the gaseous content of a patient's exhalations is measured in-situ in the patient circuit or conduit coupled to the patient's airway. In a sidestream monitoring system, on the other hand, the gas sample is transported from the patient circuit through a gas sampling line to a sensing mechanism located some distance from the main patient circuit for monitoring. As a patient's expired gases are typically fully saturated with water vapor at about 35° C., a natural consequence of the gas transport is condensation of the moisture present in the warm, moist, expired gases.\nAccurate analysis of the gaseous composition of a patient's exhalation is dependent upon a number of factors including collection of a gaseous sample that is substantially free of liquid condensate, which might distort the results of the analysis. As an expired gas sample cools during transport through the gas sampling line to the sensing mechanism in a sidestream monitoring system, the water vapor contained in the sample may condense into liquid or condensate. The liquid or condensate, if permitted to reach the sensing mechanism, can have a detrimental effect on the functioning thereof and may lead to inaccurate monitoring results. Condensed liquid in the gas sampling line may also contaminate subsequent expired gas samples by being re-entrained into such subsequent samples.\nIn addition to the condensate, it is not uncommon to have other undesirable liquids, such as blood, mucus, medications, and the like, contained in the expired gas sample. Each of these liquids, if present in the gas sample to be monitored, may render analytical results that do not accurately reflect the patient's medical status.\nThere are numerous ways in which to separate undesired liquids from the patient's expired gas stream to protect the sensing mechanism. For instance, it is known place a moisture trap between the patient and the sensing mechanism to separate moisture from the exhalation gas before it enters the sensing mechanism. The challenge, however, is to achieve the separation without affecting the characteristics of the parameters being measured, e.g., the waveform of the gas to be monitored.\nBy way of example, carbon dioxide (CO2) is effectively present only in the patient's expired gases. Therefore, the CO2 in an exhaled gas sample, transported through a gas sampling line to the sensing mechanism, fluctuates according to the CO2 present at the point at which the sample is taken. Of course the CO2 level also varies with the patient respiratory cycle. Disturbance to this fluctuation, i.e., decreases in the fidelity of the CO2 waveform, are undesirable, because such disturbances can affect the accuracy of the CO2 measurement and the graphical display of the waveform. For this reason, removal of liquids from the exhaled gas sample is desirably accomplished in such a way that it does not substantially degrade the fidelity of the CO2 waveform. Unfortunately, conventional moisture traps often disturb the waveform to a substantial degree.\nVarious other techniques have been employed to filter the expired gas stream of the undesired condensate while attempting to permit the waveform to be transported undisturbed. Such techniques include absorbents, centrifugal filters, desiccants, hydrophobic membranes and hydrophilic membranes. One well-established application which has shown some success is the use of hydrophobic hollow fibers as a filter for fluid separation. However, this application often-times still results in degradation of the waveform to some degree due to the physical requirements of the interface between the hollow fibers and the sensing mechanism.\nFurthermore, a prominent existing application of hydrophobic hollow fibers as a filter provides a disposable gas sample collection unit that connects to a reusable sensing mechanism. It this conventional filter arrangement, the gas sample collection unit, i.e., the sample cell, is physically located some distance from the filter. The gas passing through the filter is transported to the sample cell via a relatively long tube.\nThe present inventor recognized that the conventional technique of gathering and transporting the filtered gas sample to a remote sensing mechanism degrades the CO2 waveform measured at the sample cell. More specifically, the present inventor recognized that by locating the filter element some distance from the sample cell and causing the filtered gas to pass through a relatively long conduit to get from the filter to the sample cell, this arrangement effectively dampens the fidelity of the CO2 waveform, for example, by dulling the rising and falling edges of the waveform. Accordingly, a gas sampling assembly that effectively and efficiently separates moisture from a gas sample and that does not substantially degrade the waveform of expired gases measured at the sample would be advantageous."}
{"text": "Golf balls generally comprise a core surrounded by a cover and optionally intermediate layers there between. The cover forms a spherical outer surface and typically includes a plurality of dimples. The core and/or the cover may incorporate multiple layers and the core may be solid or have a fluid-filled center surrounded by windings and/or molded material. Golf ball covers may be formed from a variety of materials such as balata, polyurethane, polyurea, and/or thermoplastic compositions such as ionomer resins including SURLYN® and IOTEK®, depending upon the desired performance characteristics of the golf ball and desired properties of the cover.\nWhile conventional golf balls are white, some golfers enjoy distinguishing themselves on the course by playing a golf ball having a unique visual appearance. Accordingly, golf ball manufacturers have incorporated coloring agents such as pigments, dyes, tints, inks and the like in golf ball materials and/or coatings.\nSuch color agents typically contain chromophores. A chromophore is the group of atoms within a molecule that is responsible for the molecule's color. Specifically, bonds between the atoms allow the atoms to absorb some visible light while reflecting other visible light, with the reflected wavelengths of light attributing a certain apparent color to the molecule.\nA notable drawback with chromosphore-based coloring agents, however, is that the apparent color fades/discolors from exposure to sunlight due to photodegradation. This occurs because the atoms in a chromophore will also absorb damaging photons within the wavelengths found in sunlight (i.e., infrared radiation, visible light, and ultraviolet light), which break down the aforementioned bonds between atoms, thereby changing the atomic configuration within each chromophore. While stabilizer packages are often used to improve light stability, such improvement can be temporary. And stabilizer packages have been known to actually contribute to discoloration as well.\nAccordingly, there remains a need for golf balls having an apparent surface color that is not subject to fading or discoloration when exposed to sun light. The present invention addresses and solves this need."}
{"text": "The search for new therapeutic agents has been greatly aided in recent years by better understanding of the structure of proteins and other biomolecules associated with target diseases. One important class of these proteins is the sigma receptor, a cell surface receptor of the central nervous system (CNS) which may be related to the dysphoric, hallucinogenic and cardiac stimulant effects of opioids. From studies of the biology and function of sigma receptors, evidence has been presented that sigma receptor ligands may be useful in the treatment of psychosis and movement disorders such as dystonia and tardive dyskinesia, and motor disturbances associated with Huntington's chorea or Tourette's syndrome and in Parkinson's disease (Walker, J. M. et al, Pharmacological Reviews, 1990, 42, 355). It has been reported that the known sigma receptor ligand rimcazole clinically shows effects in the treatment of psychosis (Snyder, S. H., Largent, B. L. J. Neuropsychiatry 1989, 1, 7). The sigma binding sites have preferential affinity for the dextrorotatory isomers of certain opiate benzomorphans, such as (+)-SKF 10047, (+)-cyclazocine, and (+)-pentazocine and also for some narcoleptics such as haloperidol. The sigma receptor has at least two subtypes, which may be discriminated by stereoselective isomers of these pharmacoactive drugs. SKF 10047 has nanomolar affinity for the sigma-1 site, and has micromolar affinity for the sigma-2 site. Haloperidol has similar affinities for both subtypes. Endogenous sigma ligands are not known, although progesterone has been suggested to be one of them. Possible sigma-site-mediated drug effects include modulation of glutamate receptor function, neurotransmitter response, neuroprotection, behavior, and cognition (Quirion, R. et al. Trends Pharmacol. Sci., 1992, 13:85-86). Most studies have implied that sigma binding sites (receptors) are plasmalemmal elements of the signal transduction cascade. Drugs reported to be selective sigma ligands have been evaluated as antipsychotics (Hanner, M. et al. Proc. Natl. Acad. Sci., 1996, 93:8072-8077). The existence of sigma receptors in the CNS, immune and endocrine systems have suggested a likelihood that it may serve as link between the three systems.\nIn view of the potential therapeutic applications of agonists or antagonists of the sigma receptor, a great effort has been directed to find selective ligands. Thus, the prior art discloses different sigma receptor ligands.\nInternational Patent Application No. WO 91/09594 generically describes a broad class of sigma receptor ligands some of which are 4-phenylpiperidine, -tetrahydro-pyridine or -piperazine compounds having an optionally substituted aryl or heteroaryl, alkyl, alkenyl, alkynyl, alkoxy or alkoxyalkyl substituent on the ring N-atom. The terms aryl and heteroaryl are defined by mention of a number of such substituents.\nEuropean patent publication No. EP 0 414 289 AI generically discloses a class of 1,2,3,4-tetrahydro-spiro[naphthalene-1,4′-piperidine] and 1,4-dihydro-spiro [naphthalene-1,4′-piperidine] derivatives substituted at the piperidine N-atom with a hydrocarbon group alleged to have selective sigma receptor antagonistic activity. The term hydrocarbon, as defined in said patent, covers all possible straight chained, cyclic, heterocyclic, etc. groups. However, only compounds having benzyl, phenethyl, cycloalkylmethyl, furyl- or thienylmethyl or lower alkyl or alkenyl as the hydrocarbon substituent at the piperidine nitrogen atom are specifically disclosed. The compounds are stated to displace tritiated di-tolyl guanidine (DTG) from sigma sites with potencies better than 200 nM. As a particularly preferred compound is mentioned 1′-benzyl-1,2,3,4-tetrahydro-spiro[naphthalene-1,4′-piperidine].\nEuropean patent publication No. EP 0 445 974 A2 generically describes the corresponding spiro[indane-1,4′-piperidine] and spiro[benzocycloheptene-5,4′-piperidine] derivatives. Again the compounds are only stated to displace tritiated di-tolyl guanidine (DTG) from sigma sites with potencies better than 200 nM\nEuropean patent Application No. EP 0 431 943 A2 relates to a further extremely broad class of spiropiperidine compounds substituted at the piperidine N-atom and claimed to be useful as antiarrhythmics and for impaired cardiac pump function. The said application exemplifies several compounds, the majority of which contain an oxo and/or a sulfonylamino substituent in the spiro cyclic ring system. Of the remainder compounds, the main part has another polar substituent attached to the spiro nucleus and/or they have some polar substituents in the substituent on the piperidine N-atom. No suggestion or indication of effect of the compounds on the sigma receptor is given.\nThere is still a need to find compounds that have pharmacological activity towards the sigma receptor, being both effective and selective, and having good “drugability” properties, i.e. good pharmaceutical properties related to administration, distribution, metabolism and excretion."}
{"text": "1. Field of the Invention\nThe present invention relates to a gas laser oscillating unit, in particular, a gas laser unit provided with a gas flow passage having a junction.\n2. Description of the Related Art\nIn an axial-flow type laser oscillating unit, a gas flow passage of the unit has been improved, in which a laser medium or laser gas flows in an excitation part in order to generate a laser beam by discharging, a light reaction or chemical reaction. For example, in Japanese Unexamined Patent Publication (Kokai) No. 6-326379, a taper portion is provided on the inlet and outlet portions of a flow passage for a laser gas, in order to reduce the pressure loss of the passage. In Japanese Unexamined Patent Publication (Kokai) No. 9-199772, the shape of the front portion of an excitation part of a gas flow passage is configured in order to make the gas flow in a discharge tube a spiral flow, thereby stabilizing the discharge.\nFurther, in Japanese Unexamined Patent Publication (Kokai) No. 2003-283008 or Japanese Unexamined Patent Publication (Kokai) No. 2004-235517, two excitation parts are arranged on the same laser beam axis, and two gas flows collide with each other at a junction part between the two excitation parts. It is advantageous to merge the two flows from the two excitation parts, instead of having two separate flows, since the number of components of the flow passage may be reduced and the length of the non-excitation part may be shortened. Due to a reduction in the number of components, the cost of the laser oscillating unit may be reduced. Further, due to the shortened length of the non-excitation part, the efficiency of laser oscillation may be improved, since the energy loss of the laser beam is reduced.\nFIGS. 5 and 6 are perspective views showing an example of the structure of a junction part of a laser gas of the axial-flow type laser oscillation, as described in Japanese Unexamined Patent Publication (Kokai) No. 2003-283008 or Japanese Unexamined Patent Publication (Kokai) No. 2004-235517. As shown in FIG. 5, a tapered gas flow passages 161a is arranged between the excitation part 103a and a junction part 123, and a tapered gas flow passage 161b is arranged between the excitation part 103b and the junction part 123. The two opposing tapered passages 161a and 161b are symmetric about an axis 152 extending through a center point 151 of the junction part 123 and perpendicular to a laser axis 104. In other words, the two tapered portions are plane-symmetric about a plane including the axis 152 and perpendicular to the laser axis 104.\nTherefore, two gas flows (having the same flow rate and velocity) are generated in two flow passages, and then the two gas flows collide with each other at or near the center point 151. In this case, the state of the gas flow in the junction part 123 is unstable or easily varies. For example, the state of the gas flow at a given time t1 is represented in FIG. 5, and then the state of the gas flow at another given time t2 may be changed as shown in FIG. 6. In other words, in the state of FIG. 5, the gas flows 180a and 180b flowing from the −X excitation parts 103a and 103b are biased to the −X and +X directions, respectively, and merge in the junction part 123. On the other hand, in the state of FIG. 6, the gas flows 180a and 180b are biased to the +X and −X directions, respectively, and merge in the junction part 123. At this point, as shown in FIGS. 5 and 6, the X-direction is perpendicular to both the axis (or Y-direction) parallel to the laser axis and the longitudinal axis of a gas flow passage 162 (or Z-direction) arranged downstream relative to the junction part 123. In this example, the state of the gas flow is changed from time t1 to time t2, in other words, the state of the gas flow is unstable. Also, such an unstableness of gas flow, in two opposing gas flow passages constituting the plane-symmetrical structure, has been studied by numeric analysis, such as a finite element method of gas flow.\nThe laser gas, after flowing through the excitation part, is activated and has the property of absorbing the laser beam. Therefore, when the gas flow is unstable, the laser output and/or the shape of the laser beam mode may fluctuate. For example, in a laser machining process, the quality of a cut surface of a product may be deteriorated. Further, when the laser is used as a light source, the quantity of light may fluctuate.\nThe above Japanese Unexamined Patent Publication (Kokai) No. 6-326379 and Japanese Unexamined Patent Publication (Kokai) No. 9-199772 disclose ways to modify the gas flow passage. However, the ways are not directed to a technique for avoiding frontal collision of the gas flows as described in Japanese Unexamined Patent Publication (Kokai) No. 2003-283008 or Japanese Unexamined Patent Publication (Kokai) No. 2004-235517."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of forming a single crystalline silicon layer, a structure including the same, and a method of fabricating a thin film transistor (“TFT”) using the same. More particularly, the present invention relates to a method of forming a single crystalline silicon layer with high crystallinity, a structure including the same, and a method of fabricating a TFT using the same.\n2. Description of the Related Art\nSince poly crystalline silicon (“poly-Si”) has higher mobility than amorphous silicon (“a-Si”), poly-Si is applied not only to flat panel displays (“FPDs”) but also to various electronic devices, such as solar batteries. However, poly-Si is inferior in mobility and uniformity to single crystalline silicon.\nIn particular, single crystalline silicon is useful to a system on panel (“SOP”) structure in which a system is disposed on a display panel. The single crystalline silicon has a mobility of 300 cm2/Vs or higher. The use of single crystalline silicon with a high mobility is advantageous to formation of a high-quality switching device for a display device.\nHowever, formation of single crystalline silicon is not free from temperature limitations. That is, a process of forming single crystalline silicon cannot be performed at a temperature higher than a temperature which a base substrate, for example, a plastic substrate or a glass substrate, can resist.\nA process of forming a silicon-on-insulator (“SOI”) wafer, which is called a “smart-cut process,” includes a high-temperature annealing process that reaches a temperature of about 1000° C. Specifically, the smart-cut process includes thermally oxidizing a bare wafer with a predetermined thickness, forming a boundary layer by implanting H+ ions beneath the surface of the wafer, bonding the wafer to an additional substrate and separating the boundary layer to leave silicon on the substrate to a predetermined thickness, and performing an annealing process at a high temperature.\nIn this smart-cut process, the thermal oxidization process is performed at a temperature of 900° C. or higher, and the annealing process is performed at a temperature of up to 1100° C. Thus, there is a strong likelihood that these high-temperature processes inflict great damage on the substrate. Accordingly, the conventional method of forming an SOI wafer places a limitation on materials of the substrate and applies thermal shock even to a selected material of the substrate, thus adversely affecting the performance of a device obtained from silicon.\nAnother method of directly forming single crystalline silicon on a substrate is disclosed in “Formation of Location-controlled Crystalline Silicon” by Paul Ch. van der Wilt et al, Applied physics letters 72(12), p. 1819, 2001. This method is directed at forming single crystalline silicon on a desired location.\nSpecifically, as shown in FIG. 1, an insulating layer 2 having a hole 2a with a predetermined pattern is formed on a glass substrate (or a plastic substrate) 1, and a silicon seed layer is formed in the hole 2a. However, according to this conventional method, the surface of a silicon layer 3 is not flattened around the hole 2a. Since the silicon layer 3 has an uneven surface, it is difficult to obtain single crystalline silicon with high crystallinity."}
{"text": "The procedure for correcting the lack of homogeneity is called shimming. It consists of two steps:\n(i) mapping the magnetic field at certain points lying on the surface of a sphere centered at the magnet center, and\n(ii) correcting field inhomogeneities by driving certain currents through individual shimming coils to generate new magnetic fields that correct for the lack of homogeneity.\nShimming the NMR magnets to improve the magnetic homogeneity presently is accomplished semi-manually. More particularly, the magnetic field of the NMR imaging systems is measured for homogeneity using probes in particular regions of the static magnetic field or by imaging phantoms in particular regions. Ideally of course, the static magnetic field should be universally homogeneous within the bore of th magnet. In practice there is a small but not neglectabe lack of homogeneity in all NMR magnets. Shim coils are used to correct for this lack of homogeneity. When the shimming coil currents are applied, new mapping is done and compared. Finally the current intensities are adjusted. The measuring and shimming steps are repeated until sufficient homogeneity is obtained.\nThis is understandably a very time consuming process. NMR systems are costly, therefore, in order to be cost effective for hospitals and clinics, the systems require minimum down time and maximum operating time. Accordingly, it is in the interest of manufacturers of the NMR systems to provide for shimming the magnetic field automatically to correct inhomogeneities in as short a time as possible and with the use of the least amount of manpower as possible.\nIn the prior art, attempts have been made to use phantoms within the central region of the magnet to map the magnetic field there. Such field mapping by imaging is, for example, disclosed in a publication of the Society of Magnetic Resonance Imaging in Medicine at the second annual meeting of the society in San Francisco on Aug. 19, 1983 in a paper entitled \"Field Measurement by Fourier Imaging\" authored by A. G. Simon et al of the Columbia University of Physicians and Surgeons. In that article the authors suggest use of phantoms that comprise a single straight tube on an array of tubes in a circular disc. The article describes a phantom that is a thin water filled disc.\nThe use of the thin water filled disc measures the field homogeneity in only a single plane. A plurality of such water filled discs would test the homogeneity in a plurality of planes but would require a rather complicated and time consuming three dimensional computation to obtain a map of the magnetic field intensities. Thus the requirement for efficient inhomogeneity corrections persists."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method for automatically generating moving picture highlights of video stored in a digital video storing device using scene change detecting algorithm.\n2. Description of the Related Art\nGenerally, digital video content is delivered by a broadcasting system based on wired and wireless media such as a digital TV, and a user can store the digital video content using a storing device such as a PVR personal Video Recorder) as well as watch the digital video content.\nThough the whole content of the stored digital video may be played and watched by a user, moving picture highlights of summarized type is also provided by a program supplier or automatically generated by a system in a user side so that a user could understand the content without watching the whole video in some cases.\nThe moving picture highlights are provided for playing important part of the stored video, while representing the whole relevant video stream.\nThe moving picture highlights are provided for separately storing or playing a specific interval of the video stream, and in case that a user wants to select and watch one of many videos stored in a digital video recorder for a limited time period, a user could save time necessary for searching the desired video content by playing only moving picture highlights for each video stream.\nAlso, the moving picture highlights could also provide a preview function that may be used for a program guide apparatus necessary for a user selecting video to be recorded in the digital video storing device besides summary information of the stored video stream.\nAs the moving picture highlight should separately extract meaningful part representing the video content for a user, determination of interval for which moving picture highlights would be generated is very complicated.\nThe method of the related art adopts a method that the program supplier separately produces moving picture highlights. But, in case of a program not providing highlights, it has a problem that there is no way for a user to obtain information regarding each program within limited time if a user recorded a plurality of programs simultaneously.\nTherefore, in order to resolve the foregoing problem, a method for automatically generating these highlights is required.\nAlso, as the method of the related art generates the highlights by simply inputting time information, there is a problem that the related art method is unsatisfactory for automatic generation of the highlights such that the highlights begins at an actual scene change point.\nNamely, according to the method of the related art, it is difficult to generate moving picture highlights for the digital video in automated fashion, and extraction of the part substantially summarizing the moving picture content in viewpoint of meaning is weak and performance or reliability in automatic generation of the highlights using the extracted part is insufficient."}
{"text": "This application claims priority to DE 101 47 192.0, a German patent application filed Sep. 25, 2001, pursuant to 35 U.S.C. xc2xa7119(a)-(d).\nThis invention pertains to heat exchangers for vehicles and more particularly to such heat exchangers that include a block of heat exchange elements including flat tubes.\nFlat tube block heat exchangers have rows of flat tubes that may be alternated with corrugated fins. The ends of the flat tubes provide flow paths between tanks located at the opposite ends of the flat tubes and the flat tubes extend from one header plate to another header plate at the opposite end of the flat tubes. Typically, the flow paths created by the flat tubes extend generally perpendicular to the header plates at the ends of the tubes. The header plates are manufactured and supplied separately from the flat tubes and are typically made from relatively heavy materials. Further, the header plates occupy space at the collection tanks located at either end of the flat tube block heat exchanger.\nHeaderless flat tubes have been formed by placing two tube halves adjacent to one another where the ends of the tube halves have been deformed such that the broad sides may be connected. Drawn or welded tubes have been slit in their narrow sides in order to deform the ends and connect the broad sides. In both cases the connected broad sides form a xe2x80x9cpeak-and-valleyxe2x80x9d surface that is not conducive to directing flow from a transverse direction into the flat tubes. Examples of heat exchangers with these types of flat tubes are disclosed in German Patent Application Nos. DE 100 16 113.8 and DE 100 19 268.8, both of which are incorporated herein by reference.\nWhile these constructions can perform satisfactory for their intended purpose, there is always room for improvement. For example, the pressure loss of the medium flowing through the collection tanks and flat tubes of a block heat exchanger could be reduced in order to improve their application opportunities, especially in the vehicle field.\nIn one form, the invention provides a heat exchanger that includes a pair of collection tanks spaced opposite each other. The tanks are fluidly connected by a plurality of flat tubes that provide flow paths between the tanks. The flat tubes have opposing ends that correspond to one of the tanks and each end has two broad sides and two narrow sides. Each broad side has an inner surface and is deformed to expose the inner surface. The exposed inner surfaces of each flat tube are bonded to the exposed inner surface of any adjacent flat tube to define a fluid barrier with the corresponding tank.\nIn one form, each of the narrow sides is connected to an edge of the corresponding tank.\nAccording to one form, the deformed broad sides have two bends of about 90xc2x0 with the bonded inner surfaces arranged in parallel orientation to the broad sides of any adjacent flat tube. Alternatively, the flat tubes are arranged in oblique orientation relative to a longitudinal axis of the flat tubes and the broad sides have a first bend of about 90xc2x0 and a second bend that corresponds to the angle of the oblique orientation. In another alternative, the deformed broad sides are curved to expose the inner surface with a cross-section that is semicircular or semielliptical.\nAccording to one form, the bonded broad sides define a plurality of fluid inflow funnels to the flat tubes.\nAccording to one form, the flat tubes are formed from two half-shells.\nIn one form, each of the narrow sides includes a cut extending parallel to a longitudinal axis of the flat tubes, the cut separates the broad sides at the tube end.\nIn one form, each of the narrow sides includes two cuts that separate the broad sides at the tube ends and define a tab in each narrow side. In one form, the tab is bent inward into the flat tube.\nIn yet another form of the invention, a heat exchanger includes a pair of tanks spaced opposite each other, the tanks are fluidly connected by a plurality of spaced apart flat tubes to provide flow paths between the tanks. Each flat tube has opposing ends corresponding to one of the tanks, each end has a first broad side, a second broad side, and two narrow sides, each broad side has an inner surface and an outer surface. Each first broad side being deformed to expose the inner surface of the first broad side, and each of the exposed inner surfaces of each first broad side is bonded to the outer surface of any adjacent second broad side to define a fluid barrier with corresponding tank.\nIn another form, the invention provides a method for manufacturing a heat exchanger includes the steps of: cutting a plurality of flat tube ends to separate each tube end into a pair of broad sides; deforming each of the broad sides to expose an inner surface of the broad side; abutting the exposed inner surface of each of the broad sides to an exposed inner surface of an adjacent broad side of another tube; affixing a tank on the flat tube ends; and bonding the adjacent broad sides to each other and the flat tube ends to the tanks.\nIn one form, each flat tube end includes a pair of narrow sides in the flat tube end, and the bonding step includes bonding the narrow sides of the flat tube ends to the tank.\nIn yet another form, a method is provided for manufacturing a heat exchanger and includes the steps of: cutting a plurality of flat tube ends to separate each tube end into a pair of broad sides; deforming each of the broad sides to expose an inner surface of the broad side; abutting the exposed inner surface of one of the broad sides to an outer surface of an adjacent broad side; affixing a tank on the flat tube ends; and bonding the adjacent broad sides to each other and the flat tubes with the tanks.\nIn one form, each flat tube end includes a pair of narrow sides in the flat tube end and the bonding step includes bonding the narrow ends of the flat tube ends to the tank.\nObjects and advantages of the invention, as well as additional inventive features, will be apparent from the description of the invention provided herein and in the associated figures and appended claims."}
{"text": "The conventional carry-on case is typically a hand-carried travel case. Such cases are usually carried by a handle. It is generally necessary that this type of case be carried throughout an airport from places of departure to airplanes, from airplanes to airplanes, and from airplanes to places of arrival. For such cases there is provided a wheeled frame which is separately carried in addition to the carry-on case. The frame serves as a cart onto which the case can be strapped for transport. There are also cases that include incorporated handles and wheels so that the cases can be pulled by the handles, thereby permitting them to be towed about and transported throughout the airport. A problem with these prior art carry-on cases exists in that it is impractical, if not impossible, to use these cases as a support upon which to stack additional pieces of luggage without special devices or attachment hooks.\nThe present invention overcomes the above stated deficiencies of the prior art."}
{"text": "Heretofore, there has been known a cooling apparatus for cooling a heat exchanger with a fan driven by an engine. JP, U, 63-4400, for example, discloses a cooling apparatus comprising a heat exchanger, a propeller fan whose rotary shaft is rotated by the driving force of an engine to produce a stream of cooling air for cooling the heat exchanger, and a shroud provided downstream of the heat exchanger for introducing the cooling air to the suction side of the propeller fan, wherein a substantially disk-shaped back plate is provided just behind rotor blades of the propeller fan on the blowoff side, the back plate having almost the same diameter as an outline of the propeller fan. Such a construction is effective to avoid the occurrence of turbulence caused by interference between a main stream of the cooling air produced in the centrifugal direction on the blowoff side of the propeller fan and a reverse stream tending to return toward the heat exchanger side after being separated from the main stream, and hence to reduce noise generated by the fan."}
{"text": "1. Field of the Invention\nThe present invention relates to a printed circuit board, and more particularly, it relates to a printed circuit board having a plurality of via-holes which are provided for effectively radiating the heat generated from circuit elements which are mounted on the printed circuit board.\n2. Description of the Related Art\nRecent technology for printed circuit boards more and more tends to prepare semiconductor elements with high integration and to mount them with high density on printed circuit boards. In accordance with this tendency, it is strongly required to cool elements or to radiate heat generated from elements mounted on the printed circuit boards.\nIn actuality, a cooling apparatus is provided for a large scale electronics apparatus (for example, a large scale computer) in order to cool circuit elements. However, in the case of electronic apparatuses which are required to be small in size, it is difficult to use such a cooling apparatus because of the high cost and large size thereof. Accordingly, conventionally, in such a small scale electronics apparatus, a heat sink, for example, a metal plate, has been used for radiating the heat generated from circuit elements (i.e., heater elements) in stead of the cooling apparatus.\nIn general, the printed circuit board is made by a resin, for example, an epoxide, and this material has low thermal conductivity so that the radiation of the heat becomes worse. Accordingly, it is necessary to mount a heat sink which has a good characteristic for radiating the heat because of the high thermal conductivity thereof.\nOn the other hand, the Japanese Unexamined Patent Publication No. 62-257786 and the Japanese U.M. Publication No. 62-42273 disclose another cooling method for a purpose of low cost and small size. According to this method, the heat sink is not provided, and a plurality of through-holes (which are called \"thermally conductive through hole\" or simply \"thermal-via\") each containing a thermally conductive material, for example, a metal, are provided on the printed circuit board in order to realize good thermal conductivity in the whole printed circuit board.\nFurther, in the former document (JPP-62-257786), the through-holes are provided so as to penetrate the printed circuit board by using a drill, and copper (Cu)-plating is provided on a side surface of each through-hole. According to this structure, since the air is filled within a space provided by the through-hole, the thermal conductivity becomes extremely worse. However, since the Cu-plating is provided on the side surface of each through-hole, it is possible to improve the thermal conductivity of the printed circuit board as the whole.\nHowever, although the above method is preferable for the purpose of low cost and small size, there are some limitations, for example, the density of the through-hole and the restriction of a size of hole (i.e., diameter thereof) so that it is very difficult to obtain sufficient radiation of the heat. Accordingly, there is a problem which it is very difficult to utilize this method to elements which radiate large amount of heat.\nStill further, in the latter document (JUMP-62-42273), although this discloses the thermally conductive through-holes, it does not disclose a concrete method for improving the thermal conductivity."}
{"text": "1. Technical Field\nThe present invention relates to an image forming apparatus such as a copier, a printer, a facsimile machine, or a multi-function apparatus having one or more capabilities of the above devices, and in particular, to an image forming apparatus that can properly correct charging voltage to be applied to a charging member such as a charging roller.\n2. Background Art\nIn an electrophotographic image forming apparatus, such as a copier and a printer, it is necessary to adjust the charging voltage to be applied to the charging member in order to reliably form a quality image even though performance of the image carrier, such as a photoconductor drum, or the charging member, such as a charging roller, has been degraded over time or due to environmental changes.\nSeveral approaches have been tried. For example, the current/voltage when the rated voltage or current is applied to the charging roller is measured, and the resistance of the charging roller is obtained from the results. From a relation between the obtained resistance and the temperature detected by a temperature sensor, variation of the charged potential is forecasted and the charging voltage to be applied to the charging roller is corrected as needed. Alternatively, the temperature of the charging roller is detected and the charging voltage to be applied to the charging roller is corrected based on the detection results. Further alternatively, a surface potentiometer is used to detect the surface potential (charged potential or the exposure potential) of the photoconductor drum and image forming conditions, such as charging voltage to be applied to the charging member, are adjusted based on the surface potential of the photoconductor drum detected by the surface potentiometer."}
{"text": "This invention relates to film printers and, more particularly, to apparatus for executing light control operations during bidirectional film transport in a printer.\nFilm printers are employed to expose a reel of raw film to the frames of photographic images on a reel of master film while the master film and raw film are transported together. Light control operations are commonly executed during the exposure process in high quality film copying operations. For example, red, green, and blue color corrections are made by so-called light valves at the beginning of each scene of the master film to compensate for the different lighting conditions under which the master film was produced and different film characteristics. In order to have one scene fade out and/or the next scene fade in, a fade is executed by gradually changing the intensity of the exposing light on a frame by frame basis.\nNotches or RF strips on the film have been used for many years to mark the frames at which light control operations such as color corrections are to be executed. The color corrections are stored on a punched paper tape in the sequence in which they are to be made during film transport. Each time a notch or RF strip is sensed during film transport, the corresponding color correction is read from the tape and made by the light valves. The use of notches or RF strips to mark the master film is objectionable because the master film must be handled to so mark it, and, in the case of RF strips, the strips may come off the film.\nRecently the technique of frame count cuing has been developed to indicate the frames at which light control operations are to be executed without marking the master film. The frames of the master film are counted during transport. The frame counts at which the light control operations are to be executed are stored on the tape with the corresponding color corrections. The frame counts on the tape are compared with the actual frame count of the master film during transport to generate cue signals that execute the color corrections at the coincidence of both."}
{"text": "The invention relates to gate hinges and, more particularly, to a gate hinge having a screw for adjusting the hinged members relatively toward or away from one another.\nA number of additional features and objects will be apparent in connection with the following discussion of the preferred embodiments and examples with reference to the drawings."}
{"text": "1. Field of the Invention\nThis invention relates to a decorative lamp, and particularly to a figurative structure for clamping a decorative lamp string.\n2. Description of the Prior Art\nIn the conventional decorative lamp string for Christmas season, a plurality of sockets are mounted therein by using two or more than two power wires twisted together to connect such sockets in series; such a lamp string is subject to swinging or hanging in the air because of the sockets thereof not being fastened in place.\nIn the conventional decorative lamp string for a given festival, the figurative lamp-mounting frame is usually made of a metal, on which a plurality of socket assemblies connected in series with twisted power wires are mounted thereon. The power wires and the figurative lamp-mounting frame are usually not fastened together; as a result, the lamp string and the lamp-mounting frame are subject to separating from each other. Some of such lamp strings may be fastened in place with fastening cord; however, the sockets and the power wires are also subject to swinging and hanging in the air.\nIn the conventional decorative lamp string for Christmas season, please reference the U.S. Pat. No. 4,802,072, it is a direction fixture for decorative lamp series comprising a socket for a bulb, wires connected to said socket and a retaining ring attached to said socket, said retaining ring being provided with a notch extending longitudinally of the socket with the wires that are connected to said socket being positioned and retained in said socket so as to fixed said socket in a desired orientation and wherein said retaining ring has an outer face which is in registry with an end rim of said socket. In the aforesaid invention, the socket assemblies and the power wires are fastened in place with fastening slips, but the socket assemblies and the power wires appear to be out of order; therefore, the lamp string has to be fastened to a lamp-mounting frame with fastening cords.\nIn another conventional lamp string for Christmas season, please reference the U.S. Pat. No. 5,526,246, the front part of the figurative structure is provided with a fastening base, which is furnished with a plurality of hooks for clamping sockets respectively; then, the sockets are fastened to the figurative structure. When the bulbs wink, the figurative structure will be shown vividly. The back of the figurative structure is provided with a connection plate; as a result, such figurative structure can only be used for one-side decoration.\nIn still another conventional lamp string for Christmas season, such as U.S. Pat. No. 5,727,872, the bulb-plugging end of the socket assembly has two curved surfaces on both sides thereof, and one side thereof has two symmetrical arm plates extended out; the tail ends of such arm plates have two curved hooks respectively; the inner surface of the hooks each have a curved surface to fit to the figurative lamp-mounting frame. The outer ends of the two hooks form into an opening; each socket assembly has two arm plates to form into an opening so as to click to the metal rod of the figurative lamp-mounting frame, i.e., to have the socket assembly clamped to the metal rod of the figurative lamp-mounting frame."}
{"text": "Mooring or breasting structures are known in the maritime arts for guiding and securing vessels to a fixed location on a body of water for loading, unloading or storage. One particular type of mooring or breasting structure, known as a dolphin, consists of groups of elongate piles having a first end driven into the bed underlying the body of water and a second end extending above the body of water for contacting the vessel.\nIdeally, in order to avoid damage to either the vessel or the dolphin upon a vessel impacting the dolphin, the dolphin must possess considerable powers of resistance and also a high degree of elasticity. Oberschulte, U.S. Pat. No. 1,837,998, is illustrative of one dolphin structure which attempts to provide the combination of resistance and a high degree of elasticity. Oberschulte teaches providing a plurality of vertical hollow metal piles which are interconnected at their lower ends along a length to be imbedded into the bed underlying a body of water. The remaining portion of the piles, which extend above the bed, are unconnected along their remaining length to provide a resilient structure above the bed for elastically opposing horizontal loads. Not only does such a structure help dissipate the energy of a vessel striking fenders surrounding the piles, but by having the piles unconnected \"pull out forces\" which might otherwise lead to pulling of the piles from the bed are avoided. However, the structure taught in Oberschulte fails to optimize energy dissipation.\nThe prior art has recognized that fender piles can be utilized as part of a mooring structure in order to help dissipate the force of an impact of a vessel. Illustrative fender structures are those shown in Peterson, U.S. Pat. No. 2,420,677;\nBergfelder, U.S. Pat. No. 4,919,572; and Julian et al, U.S. Statutory Invention Registration H402.\nPeterson teaches a fender pile having a distal end rigidly embedded in a sea bed and a proximal end in abutment with a support structure such as a pier supported on a number of piles. Peterson further teaches providing shock absorbing springs at a select point along the length of the pile for allowing the pile to bend as a horizontal load is applied while preventing the pile from bending beyond its elastic limit.\nBergfelder teaches a prestressed concrete fender pile having a proximal end rigidly secured in a sea bed with a distal end abutting a pier or the like. The concrete fender pile is prestressed with high strength fiber prestressing elements which are resistant to the corrosive effects of water. Bergfelder teaches the fender pile as being able to dissipate energy by the pile bending in response to a horizontal load between the distal end and a point above the proximal end of the pile at around the level of the bed.\nJulian teaches an energy absorbing prestressed concrete fender pile which includes a rubber fender or insert between the fender pile and the pier, but which otherwise dissipates energy by deflecting in the same manner as described above with regard to Bergfelder.\nWhile each of these fender systems is somewhat effective in helping to dissipate energy which might ultimately have to be born by the pilings of a dolphin, they can result in high bending moments near the proximal end of the fender pile, which can lead to failure of the fender pile. None of these structures therefore present an optimal solution for dissipating horizontal loads.\nHolly, Jr., U.S. Pat. No. 3,852,968, teaches providing cantilevered torque arms with bumpers at their distal ends attached to vertical steel piles. Impacts on the bumpers create torsion and bending forces in the piles to help dissipate energy of an impact. However, Holly requires that the vertical piles be secured against axial rotation at their base, which can be difficult and expensive to achieve. Thus, while making effective use of torsion to absorb a load, the structure taught in Holly is prohibitively expensive in many applications.\nOne additional known way to assist in the dissipation of energy and to protect a vessel and a fender from wearing from contact therebetween is to provide a bumper on the load impacting fender. Walker, U.S. Pat. No. 3,541,800; Leblanc, U.S. Pat. No. 4,411,556; Thomerson, U.S. Pat. No. 4,338,046; Sluys, U.S. Pat. No. 4,357,891; and Aks, U.S. Pat. No. 3,486,342, teach a number of bumper structures for use on a fender.\nWhile each of the structures discussed above provides some teaching of energy dissipation for fender or dolphin structures, none of them teach an inexpensive, simple to construct dolphin structure which optimizes energy dissipation through the use of each element of the dolphin structure.\nThe present invention is directed toward overcoming one or more of the problems discussed above."}
{"text": "1. Field of Invention\nThis invention relates to a light weight light modifier having a single-piece molded housing capable of supporting several strobe-light heads, tungsten lights, HMI lights, barndoors, filter holders, diffusers, and other accessories as well as being adaptable for mounting a camera for shadow free photography.\n2. Description of Prior Art\nOne of the most important aspects of photography is the ability to control light. Any device used to control light, its intensity or direction, is called a light modifier. Those used in photography are broadly classified as follows:\na) Bounce umbrellas (modifiers) having reflective under-surface ranging from white to gold or silver metallic. The light source is reversibly mounted under the concave center of the open umbrella bouncing the light back from the reflective surface. PA1 a) Parabolic and ellipsoidal reflectors. These basic reflectors (modifiers) are usually employed with a single light source. The light bulb or a strobe light is positioned near the base of the reflector, facing the opening. The light reflects from the sides of the reflector to the area being illuminated. This type of reflector is usually made from light weight aluminum by metal spinning. Ideally the inside surface of the reflector is made of highly polished material, preferably with a textured pattern to diffuse the light and avoid hot spots. PA1 b) Soft boxes. These modifiers are made of soft, light weight fabric, supported by a framework of rods. The inside of the box is typically lined with reflective fabric and the front portion is covered with a translucent fabric. The outer-shell sides and rear portion are covered with a black fabric. The light source is mounted on the back of the box, facing the translucent front panel. The soft boxes typically employ either a single strobe light or a continuous light source, usually quartz halogen. PA1 c) Fresnel lens equipped light modifiers deliver a direct light beam, which feathers from a hard center toward a softer edge. These lights are typically made with a metal housing and have brackets for attaching accessories such as barndoors, diffusers, and filter holders. PA1 a) to enable multiple lights, either strobe based, or continuous source, to be mounted onto a single light-modifier to obtain a maximum, uniform, light from their collective output; PA1 b) to provide photographers with a light-modifier, which accepts variety of commercially available lights which, are easily mounted with adapter plates; PA1 c) to provide an economical light-modifier with a single-piece, rigid, platform capable of supporting barndoors, filters, and diffusers, and having a pivotable yoke for mounting on a stand; PA1 d) to provide a single-piece molded housing which forms a rigid platform allowing the inside of the housing to be lined with reflective panels for maximum light reflection without distorting their intended shape. PA1 e) to provide a method for mounting a camera for shadow free lighting.\nThe problem associated with the bounce umbrella is the difficulty in attaching filters for color correction or the inability to attach barndoors or other directional light control devices.\nSome of the better known umbrella designs are made and marketed by Lowell, Photogenic, Profoto, and other quality manufacturers.\nHowever, because the spinning process forces the textured surface against a tool, much of this highly reflective surface is damaged, causing less effective, uneven illumination. Typical parabolic and ellipsoidal reflectors are made and marketed by Speedotron, Norman, Profoto, and other quality manufacturers.\nThis light weight design is usually attached to a light fixture as an accessory and has no independent means for mounting on a stand. Due to its soft design, it has no hard points onto which accessories such as bardoors or filters can be readily added. Well-known soft boxes designs are made and marketed by Lowell, Photogenic, Profoto, Chimera, Westcott, and other quality manufacturers.\nWhile this light modifier successfully projects a highly controllable light beam, it is limited to a single-point light source, either a strobe or a light bulb. Only a more powerful strobe unit or a higher wattage bulb can gain additional light.\nMakers of fresnel lens equipped lights include Mole-Richardson Co., Profoto, Admagic, Sinar Bron, Elinchrom, and other quality manufacturers."}
{"text": "1. Technical Field of the Invention\nThe present invention relates generally to a sound collecting device designed to minimize electric noises caused by dust, frozen foreign substances lying on an electroacoustic transducer exposed to the air, or electromagnetic noises inputted directly to the transducer.\n2. Background Art\nFIG. 8 shows a conventional sound collecting device which consists of a horn 1 designed so as to increase in sectional area in a lengthwise direction for ease of collecting the sound wave, an electroacoustic transducer 2 (i.e., a microphone) installed in a base of the horn 1, and a preamplifier 3 connecting electrically with the transducer 2. An audio signal outputted from the transducer 2 is, as clearly shown in FIG. 9, amplified by the preamplifier 3 and outputted to an external device.\nThe transducer 2 is usually exposed to the air for catching sound waves and thus has the problems in that dust is gathered on a diaphragm of the transducer 2 with time or when the device is used in winter, it may cause the moisture in the air to be frozen solid on the diaphragm, which affects on an operation of the transducer 2, and in that since the transducer 2 needs to be exposed directly to the air, it is difficult to use a shield for protecting the transducer 2 from electromagnetic waves originating from high-voltage cables or transmission antennas, so that the electromagnetic noises are inputted directly to the transducer 2."}
{"text": "Boat ladders which are at least partly collapsible have been proposed, for example, in U.S. Pat. No. 3,590,952. One advantage of a collapsible boat ladder is that in its collapsed condition it occupies less space which facilitates storage and shipment of the ladder. It is also desirable if the length of the ladder can be adjusted. Foldable ladders have been proposed in U.S. Pat. Nos. 1,603,638, 2,758,770, and 3,286,789. However, the construction of such ladders has been unduly complicated."}
{"text": "A battery pack typically includes multiple rechargeable battery cells. High-voltage battery packs, such as those used to power a torque-generating electric machine, can generate substantial amounts of heat during sustained operation. As a result, a battery thermal management system is typically used to regulate battery cell temperature. For instance, coolant may be circulated in a closed-loop channel located near the battery cells. Thin thermal plates referred to as cooling fins may be used to help direct circulated coolant between adjacent battery cells to facilitate cooling.\nIn some battery cells, an insulating separator material may be arranged between oppositely-charged electrodes and enclosed within a sealed outer pouch filled with an electrolyte solution. The separator material, e.g., polyethylene and/or polypropylene film, helps prevent an electrical short condition while permitting the free transfer of electrical charge between electrodes. Positive and negative cell tabs of the battery cell, which are electrically connected to the respective electrodes, extend a short distance outside of the sealed pouch to form electrode terminals for the battery cell. The electrode terminals of multiple battery cells are typically ultrasonically welded together via a conductive interconnecting member positioned outside of the battery cells in order to form the battery pack."}
{"text": "1) Field of the Invention\nThe present invention relates to apparatuses and methods for preforming thermoplastic materials and, more specifically, to apparatuses and methods for bending thermoplastic sheets to form preforms for ducts.\n2) Description of Related Art\nDucts provide transport passageways for a wide variety of applications. For example, tubular ducts are widely used for air flow in aircraft environmental control systems. Similarly, ducts provide passageways for transporting gases for heating and ventilation in other vehicles and in buildings. Water distribution systems, hydraulic systems, and other fluid networks also often use ducts for fluid transport. In addition, solid materials, for example, in particulate form can be delivered through ducts. Ducts for the foregoing and other applications can be formed of metals, plastics, ceramics, composites, and other materials.\nOne conventional aircraft environmental control system utilizes a network of ducts to provide air for heating, cooling, ventilation, filtering, humidity control, and/or pressure control of the cabin. In this conventional system, the ducts are formed of a composite material that includes a thermoset matrix that impregnates, and is reinforced by, a reinforcing material such as Kevlar®, registered trademark of E. I. du Pont de Nemours and Company. The thermoset matrix is typically formed of an epoxy or polyester resin, which hardens when it is subjected to heat and pressure. Ducts formed of this composite material are generally strong and lightweight, as required in many aircraft applications. However, the manufacturing process can be complicated, lengthy, and expensive, especially for ducts that include contours or features such as beads and bells. For example, in one conventional manufacturing process, ducts are formed by forming a disposable plaster mandrel, laying plies of fabric preimpregnated with the thermoset material on the mandrel, and consolidating and curing the plies to form the duct. The tools used to mold the plaster mandrel are specially sized and shaped for creating a duct of specific dimensions, so numerous such tools must be produced and maintained for manufacturing different ducts. The plaster mandrel is formed and destroyed during the manufacture of one duct, requiring time for curing and resulting in plaster that typically must be removed or destroyed as waste. Additionally, the preimpregnated plies change shape during curing and consolidation and, therefore, typically must be trimmed after curing to achieve the desired dimensions. The jigs required for trimming and for locating the proper positions for features such as holes and spuds are also typically used for only a duct of particular dimensions, so numerous jigs are required if different ducts are to be formed. Like the rotatable tools used for forming the mandrels, the jigs require time and expense for manufacture, storage, and maintenance. Additionally, ducts formed of conventional thermoset epoxies typically do not perform well in certain flammability, smoke, and toxicity tests, and the use of such materials can be unacceptable if performance requirements are strict. Further, features such as beads typically must be post-formed, or added after the formation of the duct, requiring additional manufacture time and labor.\nAlternatively, ducts can also be formed of thermoplastic materials. A thermoplastic duct can be formed by forming a thermoplastic sheet of material, cutting the sheet to a size and configuration that corresponds to the desired shape of the duct, bending the sheet to the desired configuration of the duct, and joining longitudinal edges of the sheet to form a longitudinal joint or seam. For example, apparatuses and methods for forming thermoplastic ducts and consolidation joining of thermoplastic ducts are provided in U.S. application Ser. Nos. [. . . ] and [. . . ], titled “Thermoplastic Laminate Duct” and “Consolidation Joining of Thermoplastic Laminate Ducts,” both of which are filed concurrently herewith and the contents of which are incorporated herein by reference. Such thermoplastic ducts can be formed by retaining the thermoplastic sheet in the bent configuration until the ends are joined, and then releasing the duct so that the resulting joint continues to restrain the duct in the bent configuration. However, stresses induced in the thermoplastic material during bending can cause the duct to deform or distort from the desired configuration after joining, e.g., when released from the joining apparatus.\nThus, there exists a need for improved apparatuses and methods of preforming ducts, i.e., providing a preform configured to correspond generally to the desired configuration of the duct in a substantially unstressed condition. The method should not require the laying of individual plies on a disposable plaster mandrel. Preferably, the method should be compatible with thermoplastic ducts, including reinforced thermoplastic ducts formed from flat sheets, which provide high strength-to-weight ratios and meet strict flammability, smoke, and toxicity standards."}
{"text": "Field of the Invention\nThe present invention relates to an optical recording medium which records and reproduces information optically, and a process for producing the recording medium."}
{"text": "1. Field of the Invention\nThe present invention relates to a pneumatic gripper comprising at least one pneumatic structural element.\n2. History of the Related Art\nThe devices nearest to the present invention are known from U.S. Pat. No. 3,056,625, Timmerman (D1) and JP 05261687, Bridgestone (D2).\nD1 describes a gripper for goods which is configured as a clamp and has grippers held movably on its lateral vertically disposed sections by means of hinges. Located between each vertically disposed section and the associated gripper is an inflatable bellows which, under pressure, pushes the gripper away from the vertical section towards the inside of the gripper so that the goods are grasped by the gripper.\nD2 also describes a gripper configured in the manner of a clamp with lateral vertically disposed sections. A supporting arm runs parallel to each of the vertical sections, on the inner side thereof, said supporting arm for its part being movably hinged in the cross member of the gripper by means of a hinge. Inflatable cushions are provided on each supporting arm which, when filled with compressed air, press the supporting arms away towards the inside (whereby said cushions also move inwards since they are disposed on the supporting arms).\nThese known grippers have the disadvantage of a rigid structure, for example, with the consequence that they can only be used in a vertical position."}
{"text": "An abrasive disk is usually prepared by coating an adhesive on a back pad disk, and bonding the back pad to a disk form of a coated abrasive body (comprised of a backsheet and a layer of an abrasive material), followed by heat-pressing. The back pad is generally prepared using glass fibers for dimensional stability, and a conventional abrasive disk comprising a glass fiber textile-containing back pad is shown in FIG. 1.\nHowever, the glass fiber textile has problems in that it is heavy, expensive and stiff, which limits the use of such an abrasive disk.\nIt is also known that abrasive materials in the edge of the abrasive disk wear down quicker than those in other part of the disk during polishing, leading to lowering of the abrasion efficiency. Thus, the worn abrasive material region is usually ground out together with the unused abrasive materials on the other part of the abrasive disk, by a procedure known as “dressing”. This dressing operation is generally performed in several steps, during which the glass fiber textile generates a glass-fiber dust which irritates the skin and respiratory system of the worker. Further, the glass fiber textile has unsatisfactory wear resistance, which leads to a poor productivity and an increase in the manufacturing cost."}
{"text": "The present specification relates generally to the field of displays. More specifically, the specification relates to virtual displays.\nVirtual displays can provide information that is viewable in virtual space for a user of equipment, such as aircraft, ships, boats, naval craft, medical equipment, robotic equipment, remote vehicles, unmanned vehicle systems (UVSs), training simulators, entertainment systems, military equipment, land vehicles, etc. The information can include navigation parameters, guidance parameters, equipment parameters, location information, video information, remote views, symbology, etc.\nVirtual displays can be Near Eye Displays (NEDs), such as Head Mounted Displays (HMDs) (e.g., head worn displays, helmet mounted displays and head worn displays) or Head Up Displays (HUDs) with a fixed combiner near the eye position. Virtual displays can be utilized to provide images to an operator or user (e.g., a pilot in a cockpit). In aircraft applications, HUDs generally include a fixed combiner, an optical projector, an image source, and a HUD computer. HMDs generally include a head worn or helmet mounted combiner, optical projection elements, an image source, a HMD computer, and a head orientation sensor. The HUD or HMD computer causes the image source to provide an image which is projected to a combiner. The combiner provides a collimated image to the pilot. The image can include enhanced vision images, flight symbology, targeting data, flight instrument data, synthetic vision images, head up display (HUD) data, etc.\nCockpit and other display technologies have utilized non-virtual displays such as gauges and panel displays (e.g., head down displays (HDDs) in the cockpit environment). The non-virtual display technology has migrated from a multiplicity of independent gauges to a few large panel, non-virtual displays (e.g., large format HDDs in the cockpit environment). The large format HDDs can represent and concentrate information that used to be apportioned to different gauges and smaller HDDs. While this display technology allows for denser and more flexible display of multiple information streams, the denser, larger display formats can present several drawbacks. First, as information density on each panel is increased, the failure of a single display panel can cause degradation in cockpit workflow, pilot workload, and the amount of information provided to the pilot. Second, the larger, denser displays cannot easily direct attention to particular warnings on the HDDs and/or to locations outside of the cockpit and/or off the HDDs. Third, larger displays cannot be designed to cover the entire cockpit area. Dead space or unused areas in the cockpit cannot be filled in with display information due to shape and size constraints. Large HDDs often include bezels which take up space in the cockpit and cannot display information for the pilot. Bezels associated with conventional HDDs can prevent a seamless display experience.\nThus, there is a need for a low cost, lightweight virtual display system for use with a heads down display (HDD). There is also a need for a virtual display system that provides a seamless display system. There is further a need for near eye display system and method that can be easily integrated in the design of a cockpit without requiring extra display space. There is further a need for a near eye display that can provide display redundancy in the event of a malfunction. There is also a need for a virtual display system and method that is optimized to direct attention to warnings in the cockpit or to locations outside of the cockpit. Yet further, there is a need for a near eye display system that displays types of information in positions that are appropriate for the particular type of information.\nAccordingly, it would be desirable to provide a display system and/or method that provide(s) one or more of these or other advantageous features. Other features or advantages will be made apparent in the present specification. The teachings disclosed extend to those embodiments which fall within the scope of the appended claims, regardless of whether they accomplish one or more of the aforementioned advantages or features."}
{"text": "The present disclosure is related to cache memory, and more particularly, to a prefetching data to a lower level cache memory address translation."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of analyzing wear particles in used lubricants. More specifically, the present invention discloses a sample preparation system for separating wear particles from such fluids by size and magnetic characteristics.\n2. Statement of the Problem\nThe present invention is designed to filter and separate wear particles suspended in used lubricants or hydraulic fluids by particle size and magnetic characteristics. These samples can then be analyzed qualitatively or quantitatively using a conventional energy dispersive x-ray fluorescence (EDXRF) system or other analysis techniques.\nMore specifically, the present invention is designed to meet the particular needs of the Joint Oil Analysis Program (JOAP) laboratories. These laboratories belong to the Army, Navy, Air Force, and other Department of Defense (DoD) agencies. All of these laboratories are part of an interagency cooperative effort to implement an effective condition monitoring response to threats, and to increase safety of service personnel and the longevity of the hardware transporting the personnel. It is the responsibility of the JOAP Technical Support Center (TSC) to set the equipment standards for analysis for each of the individual condition monitoring and oil analysis programs operated by the Army, Navy, Air Force, and other DoD agencies. If a piece of equipment is to be used in the condition monitoring program of any DoD agency, the analysis system or technique to be used must first receive approval from the JOAP-TSC in Pensacola, Fla.\nCurrent atomic emission spectrographs used by the JOAP are most sensitive to particles having a size of approximately 10 microns or less. Particles larger than 8 to 10 microns in size cannot be vaporized by the high-voltage arc of the atomic emission spectrograph and therefore are not part of the sample analysis. A growing number of people in the engine condition monitoring field are now of the opinion that analysis of a wide range of particle sizes, including larger particles of 10 microns and up, is far more indicative of abnormal wear and provides the best indication of impending, possibly catastrophic failure. This seems reasonable because small metal particles present in the oil are the result of both normal wear and large particles being ground into small particles by the mechanism. Therefore, any analytical technique that is capable of analyzing only small particles is going to be less effective in predicting a need for engine maintenance, and will only occasionally be able to predict the impending catastrophic failures that are most hazardous. It is estimated that current atomic emission spectroscopy techniques detect abnormal wear only an average of once in every 5,000 analyses performed on used oil samples. There have also been several instances where an aircraft has had an in-flight failure shortly after having an oil analysis by atomic emission spectroscopy indicating no unusual quantities of wear metal. Visual inspection of the oil filter and other interior portions of the engines after the crash revealed numerous large particles, chips, and chunks of metal debris. An atomic emission spectroscopy analysis of the used oil collected after the crash still showed nothing abnormal.\nAn analytical technique that is sensitive to a wide range of particle sizes, including large particles, is advantageous and would allow detection of impending failures that are not observed with the current technique. This translates to a reduction in maintenance costs, increased safety for personnel, and fewer in-flight failures. EDXRF is sensitive to all particle sizes and therefore offers an opportunity to significantly improve the early detection of abnormal wear in aircraft. As a result of these concerns, there is a need for a system to prepare filtration samples in which particles are separated into multiple categories by size and magnetic characteristics. In addition, the system should allow the operator to easily identify any samples containing significant quantities of oversize particles.\nAlthough the present invention was specifically developed to support aircraft engines used for helicopters and jet fighters, it should be expressly understood that the invention is also applicable to analysis of used lubricants and hydraulic fluids from commercial aircraft, ground-based equipment such as heavy construction equipment, trucks, power generation stations, ocean liners, other types of ships, and high performance automobiles.\nA number of particle filtration systems have been used in the past for analysis of used oil and lubricants, including the following:\n______________________________________ Inventor Patent No. Issue Date ______________________________________ Thornton et al. 4,555,331 Nov. 26, 1985 Cox et al. 4,550,591 Nov. 5, 1985 Luria 4,169,677 Oct. 2, 1979 Westcott 4,047,814 Sep. 13, 1977 Vobach et al. 2,105,851 Jan. 18, 1938 ______________________________________\nThornton et al. disclose a semi-automatic quantitative filtration assembly having means to measure and present a known quantity of fluid for filtering. An in-process fluid holding tank (e.g., moat) adjacent to the filter medium receives any excess fluid during filtering to prevent excess fluid from intermixing with the known quantity of fluid to be filtered.\nCox et al. disclose an apparatus for monitoring particulates in a fluid. The fluid flows through a filter 28. A pressure sensing means 46 senses the fluid pressure difference across the filter. A processor evaluates the rate of change of the pressure difference in order to give an indication of the particulate matter levels in the fluid.\nLuria discloses a system for analyzing used oil in which a plurality of samples are separated according to the different mobility rates of its particles through a liquid medium. The samples are separated by passing them through a large mesh screen at the bottom end of a tubular holder in which the centrifuging is effected.\nWestcott discloses a system for detecting and analyzing particulate matter suspended in a fluid. The fluid flows over a collecting substrate in the presence of a magnetic or electric field having an intense gradient. The particles are deposited on the substrate in accordance with their size with the larger particles being deposited first and the smaller particles being deposited last.\nVobach et al. disclose a system for measuring contaminants in lubricants. A magnetic field is used to separate magnetic metal particles suspended in the lubricant.\n3. Solution to the Problem\nNone of the prior art references uncovered in the search show a sample preparation system for analysis of wear particles in lubricants that produces filtration samples in which wear particles are separated into multiple categories by size and magnetic characteristics. In addition, the system allows the operator to easily identify any samples containing significant quantities of oversize wear particles."}
{"text": "This application is a continuation of U.S. patent application Ser. No. 09/859,296, filed May 17, 2001, entitled “DOMESTIC ORIGINATION TO INTERNATIONAL TERMINATION COUNTRY SET LOGIC”, now U.S. Pat. No. 6,618,475, which has a common assignee and inventorship to the present application, and the above noted patent application, is incorporated by reference in its entirety.\n1. Technological Field\nThe present application relates generally to fraud control in telecommunications systems and, in particular, to preventing fraud in calls from a domestic origin point to an international terminating point in a long distance telecommunications network.\n2. Description of the Related Art\nThe telecommunications industry has experienced significant changes in the way that customers are billed for their telephone calls. From the once simple method of billing the originating caller, many methods have been developed, allowing greater flexibility for the telecommunications customer. A predominant method for making telephone calls away from home or the office is by utilizing the telephone calling card to charge the call.\nCalling card customers may use any telephone facility, including public facilities, to make a call that will be charged to their account. The process of making calls using a calling card typically includes dialing an “800” number, waiting for an audio prompt, and then entering an account number and a Personal Identification Number (PIN) into a telephone key pad device. The “800” (and now “888”) number phone calls are one type of a category of phone calls called “special service” calls. These special service calls, which include “700”, “800/888”, and “900” number calls, allow contemporary telecommunications networks to provide many services beyond direct distance dialing. It is the long distance carriers that provide this special service call processing, which allows for toll-free calls, calling card calls, special rate calls, etc.\nFollowing the example of a calling card call, once the account number and PIN have been entered, the calling card customer can make one or more calls from whatever location the customer is dialing in from. These calls are subsequently charged to the customer's calling card account. Calling cards can also be used to avoid having to pay additional surcharges when making calls from certain public facilities such as hotels and telephone booths.\nAs with many new technologies, the ease and flexibility of the use of calling cards has led to abuse, and has consequently brought about new types of fraud. Calling card fraud costs businesses (and consumers) millions of dollars annually. Current security mechanisms, while effective, are not fail-safe, and protection mechanisms for consumers and businesses require improvement to stem these fraud-related losses.\nThere is a virtual underground industry in stolen calling cards and authorization codes. The multitude of ways that calling cards and authorization numbers find their way into unscrupulous hands need not be discussed here, but suffice it to say there is no end to the ingenuity of the criminal mind. One example of calling card fraud is the technique of “surfing” banks of public telephones, such as are at airports. Criminals “surf” by looking over the shoulders of legitimate card users as they key in the account number and PIN. Then they sell or distribute these numbers and rampant fraud results. In some cases, a single account may incur charges in excess of $100,000 in a single weekend. Calling card fraud and other forms of fraudulent use present pervasive problems for telephone carriers, particularly long distance carriers.\nOne method of fraud control is to simply remove calling card numbers against which it is suspected that fraudulent calls are being charged. In order to recognize fraudulent calls, a “billing number”—a billing product and an account number, such as a calling card, pre-paid phone card, etc.—is monitored over time. For example, where the number of domestic calls placed within a certain amount of time using the same billing number exceeds a certain threshold, an alert is generated. International calls may have a lower threshold so that fewer calls within the time period generate an alert. In addition, the threshold may be further adjusted for calls to countries where a high percentage of fraudulent calls are directed.\nAnother method of fraud control is to identify particular origin points that are linked to suspicious activity and to block certain calls from those particular origin points. For example, a large number of long duration calls to China may be generated from an exchange in Manhattan. This would generate a threshold alert, which is typically sent to a fraud analyst. A fraud analyst would be stationed at a fraud control console 100, as shown in FIG. 1. The fraud analyst analyzes the alert and the history of that exchange in order to determine whether or not to block that exchange from calling China. If the fraud analyst decides that there is fraudulent activity, he sets up a block on that exchange which will prevent subsequent calls to China or other international destinations that the fraud analyst selects.\nThe present invention concerns this type of blockage and, in particular, blocks on special service calls that originate domestically and terminate internationally. An example of this type of special service call is shown with reference to FIG. 1. The caller, using telephone 111, makes a calling card call by dialing a number in the format of 1-NPA-NXX-XXXX. NPA stands for Number Plan Area, often referred to as the “area code”, which defines the geographic region of the number; NXX is the terminating exchange, typically identifying a switch within the geographic region; and XXXX is the unique station designation. For most calling cards, the number will take the form 1-800-NXX-XXXX, where the “800” signifies that the call is a special service call, rather than a geographic region. The call is routed through Local Exchange Carrier (LEC) 120. LEC refers to local telephone companies, such as the Regional Bell Operating Companies (RBOCs), which provide local transmission services for their customers. Because of the 1-800 format of the dialed number, the routers in the LEC will forward the call to the network of the appropriate long distance carrier (or Inter-Exchange Carrier IXC) 130. Special service telephone calls, such as “800” number calls, are provided by IXCs, such as MCI-Worldcom.\nReturning to our call, after switching through LEC switches 122 and 124, the “800” number is routed from POP (Point-of-Presence) switch 125 into the IXC 130, and then through IXC switches 131 and 136, to a bridge switch 135. The purpose of the bridge switch 135 is to receive calls from the IXC network and bridge them to the Automatic Call Distributor (ACD) 140 and, ultimately, into the Intelligent Services Network platform (ISN) 150. Because special service calls require special call processing, they are typically routed to a call processing platform, such as the ISN platform 150. There are a number of ISNs within the IXC, but, for the purpose of understanding the present invention, one ISN will suffice.\nAn exemplary and simplified diagram of the ISN platform 150 will now be described with reference to FIG. 2. The ACD 140 is under the direct control of the Application Processor APP 156, which is a general purpose computer that functions as the central point for call routing control in the ISN 150. When the “800” number call arrives at the ACD 140, the ACD 140 makes a request to the APP 156 for directions as to how the call should be handled. Such a request would usually be accompanied by information concerning the call; i.e. the Automatic Number Identifier (ANI) of the caller and the destination number of the call. The APP 156 would recognize by the “800” prefix of the destination number that the call is a special services call and, consequently, the APP 156 would instruct the ACD 140 to deliver the call to the appropriate queue. In this case, assuming that the call is to a calling card “800” number, the call would queue up to the Automatic Response Unit (ARU) 152. The ARU 152 comprises two components, one to process the call, the other to prompt the caller with a voice response system. It is the ARU 152 that will ask the caller for the required final destination number, calling card number, and PIN. When a live operator is required, the call is routed to the Manual Telecommunications Operator Console (MTOC) 154. Whether the call is routed to the ARU 152 or the MTOC 154, the same informational decisions will have to be made. In other words, regardless of whether it is entered by the operator at the MTOC 154 or by the caller at her telephone 111 to the ARU 152, items such as the calling card account number will have to be entered.\nDuring the course of servicing a call, the need often arises to “park” a call on the ACD 140. When a call is parked on the ACD 140, the call is active, i.e., there is a party on the call with an established voice channel connected to the ACD 140. The call is monitored and maintained at the ACD. Once a call is parked at the ACD, it is no longer under direct control of either the ARU 152 or the MTOC 154 that parked the call. This allows the facilities at the ISN 150 to be freed up to perform other tasks or services. When call processing is completed, and the call is authorized and validated, the call is released from the ACD 140 and the bridge switch 135 to the automated switching of the IXC network 130. As shown on FIG. 1, the call is then connected through IXC switch 137 to a telephone 199 in China.\nNow, a simplified and exemplary call processing procedure will be described with reference to FIGS. 2, 3A, and 3B. Many steps that are required for call processing have been eliminated from the description as unnecessary for the understanding of Country Set Logic. Assuming that the special services call is a calling card call from a domestic origin to an international destination, the caller needs to enter her account number, PIN, and the terminating ANI. It is assumed that all of this data is input before the procedure begins, but, as one skilled in the relevant art would know, some of the data could just as well be entered during the procedure. Following this example, once input is complete, the access code is looked up in an access-level database, such as the Authorization Property Database (AUTH PROP) 168, at step 300 in FIG. 3A. The access code is the original 1-800-NXX-XXXX dialed in to access special services, and an access-level database is a database keyed to the various access codes.\nRecords in the Authorization Property Database 168 contain various items keyed to the access code, including operator scripts, billing products, and options. The AUTH PROP records also contain a field for Country Set Logic in order to indicate limitations on international destinations. Basically, Country Set Logic consists of this extra field where a term in the form CSETX (where X represents a number from 1 to 999) can be placed. For example, the access code “1-800-555-6543” might contain “CSET16” in its Country Set field (CSET). This means that, when the destination number is looked up in the international database, if the term “CSET16” appears in the international database, the call will be blocked. In step 305 of FIG. 3A, it is determined whether there is an access-level CSET term in the AUTH PROP 168. If there is, the CSET term is saved in step 307. If either there is no CSET term in step 305, or after the CSET term is stored in step 307, the call processing continues at step 310.\nAt step 310, the billing number associated with the customer account is looked up in the Billed Number Screening (BNS) database 160. The BNS contains records keyed by billing numbers and has flags to indicate various limitations on particular billing numbers. It is determined whether the billing number is flagged in step 315. If the billing number is flagged, the call may be re-routed to an MTOC 154, a fraud analyst at a fraud console 100, or simply disconnected. If the billing number is not flagged in step 315, the exchange of the originating ANI (the prefix NPA-NXX of the originating number) is looked up in an exchange-level database, such as the Exchange Master (X-MASTER) database 162, in step 320. The X-MASTER has records keyed on the various NPA-NXXs and also includes flags that indicate various limitations on particular exchanges. The records in the X-MASTER also contain a CSET field in order to indicate limitations on international destinations. In the same manner as the access-level database, if the originating exchange NPA-NXX contains “CSET32” in its CSET field and the term “CSET32” appears in the international database, the call will be blocked.\nThe CSET logic is the primary focus for the rest of this application. With this in mind, the X-MASTER 162 and other databases discussed here would likely be accessed for other purposes, such as viewing other flags and fields. For example, X-MASTER 162, like the BNS 160 in step 310, is typically accessed to determine if there are other blocking flags on the originating exchange. This step, and others, have been left out as extraneous to an understanding of Country Set logic and the present invention.\nReturning to the call processing procedure, in step 325, it is determined whether there is a CSET term in the particular NPA-NXX record in X-MASTER 162. If there is, the CSET term is saved in step 327. In step 327, if there was a previously stored access-level CSET, it is deleted and replaced with the exchange-level CSET. This establishes greater granularity, because the exchange-level is much narrower than the access-level. If either there is no CSET term in step 325, or after the CSET term is stored in step 327, the call processing continues at step 330.\nIn step 330, the originating ANI is looked up in an ANI-level database, such as the ANI Property database (ANI PROP) 169. The ANI PROP 169 contains records keyed to ANIs, and the records contain flags, fields, and other information unique to that ANI. This provides the greatest granularity, because a particular payphone can be blocked using an ANI-level database. There is also a CSET field in the ANI PROP 169 records and, in step 335, it is determined whether there is a CSET term in the originating ANI's record in the ANI PROP 169. If there is, the CSET term is saved in step 337. In step 337, if there was a previously stored access- or exchange-level CSET, it is deleted and replaced with the ANI-level CSET.\nSince the focus is on the CSET logic, if there is either no CSET term in step 335, or after the CSET term is stored in step 337, the procedure jumps to step 350 on FIG. 3B, leaving out many call processing details. In step 350, the destination number is checked against the International Country Code Database (INTERNAT'L COUNTRY) 164. The INTERNAT'L COUNTRY has records keyed on the various international country codes and also includes flags that indicate various limitations on the particular countries. The records in the INTERNAT'L COUNTRY also contain a field for Country Set Logic in order to indicate limitations on international destinations. In step 355, it is determined whether there is a matching CSET in the INTERNAT'L COUNTRY. If there is a matching CSET in step 355, the call is blocked. If not, call processing continues until completion. During the continuation of call processing, other databases, including the International City Code (INTERNAT'L CITY) database 166, are accessed.\nThis method is effective in eliminating fraudulent calls made from origin points that have been recognized as generating a large amount of fraudulent calls to particular international destinations. Typically, fraud control 100 maintains a fraud-to-revenue ratio in relation to particular exchanges calling particular countries. Once this fraud-to-revenue ratio reaches a certain threshold, some form of CSET logic is placed on the originating exchange/destination country combination. This type of block makes sense because most hackers will move from phone to phone within a certain area. Thus, the conventional method eliminates a great deal of fraud; however, it is troublesome to legitimate callers within that exchange. A calling card customer making a non-fraudulent call within that exchange will be blocked from calling that international destination. In addition, blocking the entire country will sometimes cast a much larger net than is needed for the task.\nTherefore, there is a need to allow legitimate callers to make calls from blocked exchanges to international destinations. In addition, there is a need to permit finer granularity in blocking calls to international destinations."}
{"text": "Embodiments presented herein are related to electronic devices, and more specifically, to determining whether a speech user interface would be a constructive addition to a graphical user interface (GUI) of an application.\nA GUI, is a type of user interface that allows users to interact with electronic devices through graphical icons, visual indicators, or the like, which are collectively referred to herein as GUI objects. Users may interact with the electronic device through direct manipulation of the GUI objects displayed in the GUI. Beyond computers, GUIs are used in many handheld mobile devices such as smartphones, MP3 players, portable media players, gaming devices, tablets, and smaller household, office and industrial devices.\nA voice user interface (VUI) makes user interaction with the electronic device possible by using voice/speech. In general, the more complex the user's speech interactions, the more challenging it is for the VUI to recognize those speech interactions. On the other hand, a VUI may be beneficial for handling quick and routine user interactions with the electronic device.\nSome electronic devices, such as mobile electronic devices, are ripe for speech based interfaces due to the typically small display size. In such devices, challenges exist both in providing input (e.g. via small virtual keyboards) and output (e.g. must tap/swipe/pinch to scroll through content). Additionally, speech can provide a hands-free mechanism to interact with the GUI of an application which can be useful for the user. Thus, it may be beneficial to determine whether the user's efficiency and/or experience of interacting with the GUI would be increased by adding a VUI."}
{"text": "Typically, subsea equipment used in oil and gas applications must be lowered to a wellhead, a subsea equipment or system, such as a Christmas tree, or other site at the seabed. One type of subsea equipment that is lowered into the sea for installation may be a flow control module, for example. A flow control module is typically a preassembled package that may include a flow control valve and a production fluid connection that can mate with a hub on a subsea equipment or system, such as a Christmas tree. The hub on the Christmas tree may include a production fluid conduit to allow for the flow of production fluid from the well. The Christmas tree is typically mounted to a wellhead.\nTypically, the flow control module may also include electrical and hydraulic connections as well as gaskets. The electrical and hydraulic connections may be used to control and serve components on the tree, such as valves. These connections or gaskets may be assembled on a flange of the production fluid connection for mating with corresponding connections on the tree hub. A stab and funnel system between the tree and flow package is typically used to align the production conduit and the several connections on the flow control package with those on the tree hub. Hard landing the flow control package on the tree may damage the connections at the hub, given the heavy weight of many equipment packages. To reduce the possibility of damage to the connections, the flow control module can be soft landed onto the tree. Soft landing is carried out by a running tool having a complex system of hydraulic cylinders and valves that slow the descent of the flow module package as it is landed onto the tree. However, the use of such soft landing running tools can be very expensive.\nA need exists for a technique to achieve soft landing of subsea equipment without the use of a running tool."}
{"text": "This invention concerns an attachment device for a woodworking plane, and is further directed to a woodworking plane provided with an attachment which enables said plane to produce straight edge surfaces of various desired angles with respect to the face of a workpiece.\nIn the course of woodworking operations, a manually operated plane is frequently utilized to smooth or straighten surfaces, particularly the end or edge surfaces of a relatively large flat workpiece having at least one substantially straight edge. In the case of panels such as doors, and in other woodworking applications, it is often found desirable to cause the edge surface to be angled at other than 90.degree. with respect to the front or rear surfaces of the panel. It is sometimes sought to cause said edge surfaces to deviate by about 5 to 20 degrees from the 90 degree configuration initially present at the edges of the panel.\nNumerous devices have been proposed as attachments to woodworking planes to facilitate the cutting of edge surfaces at accurate angles with respect to face surfaces. Such devices, however, have either been inaccurate, expensive or cumbersome to use.\nAccordingly, it is an object of this invention to provide an attachment device which will enable a woodworking plane to cut straight edge surfaces of a workpiece at an accurate and adjustable angle with respect to a face surface of said workpiece.\nIt is another object to provide a device of the aforesaid nature of simple and rugged construction which may be economically manufactured.\nIt is a further object of the present invention to provide a device of the aforesaid nature which may be rapidly applied to and removed from a plane, and is simple to operate.\nIt is a still further object of this invention to provide a woodworking plane equipped with an improved attachment device which enables said plane to cut straight edge surfaces of a workpiece at an accurate angle with respect to a face surface of said workpiece.\nThese objects and other objects and advantages of the invention will be apparent from the following description."}
{"text": "The malfunctioning of protein kinases (PKs) is the hallmark of numerous diseases. A large share of the oncogenes and proto-oncogenes involved in human cancers encode for PKs. The enhanced activities of PKs are also implicated in many non-malignant diseases, such as benign prostate hyperplasia, familial adenomatosis, polyposis, neuro-fibromatosis, psoriasis, vascular smooth cell proliferation associated with atherosclerosis, pulmonary fibrosis, arthritis glomerulonephritis and post-surgical stenosis and restenosis.\nPKs are also implicated in inflammatory conditions and in the multiplication of viruses and parasites. PKs may also play a major role in the pathogenesis and development of neurodegenerative disorders.\nFor a general reference to PKs malfunctioning or deregulation see, for instance, Current Opinion in Chemical Biology 1999, 3:459-465.\nNeuroblastoma, a pediatric malignancy of the sympathetic nervous system, is characterized by clinical and biological heterogeneity. Approximately one-half of neuroblastoma patients present with advanced-stage disease, and despite intensive multimodality therapy, including myeloablative regimens, survival for these children is less than 40%. Identification of tumor targets and advances in target-specific therapies with minimal non-specific toxicity are needed for this patient population. The Trk family of receptor tyrosine kinases is critical for neuronal survival and differentiation during the development of the nervous system. The Trk receptors are differentially expressed in human neuroblastoma and likely play a central role in tumorigenesis and/or cell survival. TrkA is highly expressed by neuroblastomas with favorable biological and clinical features, and expression is associated with patient outcome. In contrast, TrkB expression is restricted to a malignant subset of neuroblastomas. Co-expression of TrkB and its ligand, BDNF, in the majority of neuroblastomas, provides a potential autocrine survival pathway in biologically aggressive, high-risk tumors. Additionally, the recent identification of a constitutively active TrkA splice variant (TrkAIII) that is preferentially expressed in advanced-staged tumors highlights the complex role of Trk signaling in neuroblastoma biology and its potential as a therapeutic target.\nNeurotrophin signaling through the Trk family of receptor tyrosine kinases (RTKs) plays a critical role in the development, maintenance and function of the nervous system. Activation of these receptors regulates cell survival, proliferation, migration, differentiation, and apoptosis during development. They exert this influence by modulating the responses of neurons to the neurotrophin family of growth factors in a temporally and spatially regulated manner. The neurotrophins nerve growth factor (NGF), brain-derived neurotrophic factor (BDNF), and neurotrophin-3 (NT3) are the cognate ligands for TrkA (NTRK1), TrkB (NTRK2), and TrkC (NTRK3), respectively.\nNeuroblastoma, a common pediatric tumor of the postganglionic sympathetic nervous system, provides an ideal model for the study of Trk signaling and inhibition in cancer. Neuroblastomas are characterized by clinical heterogeneity, from spontaneous regression in infants to relentless progression in older children. The prognosis for these latter patients remains poor, with three-year event-free survival (EFS) probabilities of 30-40% (5-7). Indeed, neuroblastomas can be classified into distinct subsets based on genetic alterations and biologic features (8), and the expression of Trk receptors likely contributes to these distinct behaviors.\nExpression of TrkA in neuroblastoma cell lines has been shown to mediate neuronal differentiation, growth arrest and inhibition of angiogenesis in response to NGF. In contrast, unfavorable neuroblastomas frequently express TrkB and its ligand BDNF, which together comprise an autocrine or paracrine survival pathway. These tumors typically have gross segmental chromosomal aberrations including amplification of the MYCN proto-oncogene. The TrkB/BNDF pathway promotes cell survival, protects cells from injury, and blocks chemotherapy-mediated cell death in vitro. Although a number of genes are likely involved in the development and clinical behavior of favorable and unfavorable neuroblastomas, the pattern of Trk gene expression (TrkA versus TrkB) likely plays a role.\nRecent literature has also shown that overexpression, activation, amplification and/or mutation of Trk's are associated with many cancers including neuroblastoma (Brodeur, G. M., Nat. Rev. Cancer 2003, 3, 203-216), ovarian cancer (Davidson. B., et al., Clin. Cancer Res. 2003, 9, 2248-2259), breast cancer (Kruettgen et al, Brain Pathology 2006, 16: 304-310), prostate cancer (Dionne et al, Clin. Cancer Res. 1998, 4(8): 1887-1898), pancreatic cancer (Dang et al, Journal of Gastroenterology and Hepatology 2006, 21(5): 850-858), multiple myeloma (Hu et al, Cancer Genetics and Cytogenetics 2007, 178: 1-10), astrocytoma and medulloblastoma (Kruettgen et al, Brain Pathology 2006, 16: 304-310) glioma (Hansen et al, Journal of Neurochemistry 2007, 103: 259-275), melanoma (Truzzi et al, Journal of Investigative Dermatology 2008, 128(8): 2031-2040, thyroid carcinoma (Brzezianska et al, Neuroendocrinology Letters 2007, 28(3), 221-229.), lung adenocarcinoma (Perez-Pinera et al, Molecular and Cellular Biochemistry 2007, 295(1&2), 19-26), large cell neuroendocrine tumors (Marchetti et al, Human Mutation 2008, 29(5), 609-616), and colorectal cancer (Bardelli, A., Science 2003, 300, 949). In preclinical models of cancer, Trk inhibitors are efficacious in both inhibiting tumor growth and stopping tumor metastasis. In particular, non-selective small molecule inhibitors of Trk A, B and C and Trk/Fc chimeras were efficacious in both inhibiting tumor growth and stopping tumor metastasis (Nakagawara, A. (2001) Cancer Letters 169:107-114; Meyer, J. et al. (2007) Leukemia, 1-10; Pierottia, M. A. and Greco A., (2006) Cancer Letters 232:90-98; Eric Adriaenssens, E. et al. Cancer Res (2008) 68:(2) 346-351) (Truzzi et al, Journal of Investigative Dermatology 2008, 128(8): 2031-2040. Therefore, an inhibitor of the Trk family of kinases is expected to have utility in the treatment of cancer.\nVarious gene rearrangements of the Trk gene have been implicated in human malignancies. For example, the MPRIP-NTRK1 and CD74-NTRK1 gene rearrangements have been implicated in the development of non-small cell lung cancer. Gene rearrangements TPM3-NTRK1, TGF-NTRK1 and TPR-NTRK1 have been implicated in the development of papillary thyroid cancer. The TPM3-NTRK1 gene rearrangement has been implicated in the development of colorectal cancer. NTRK1, NTRK2 or NTRK3 gene rearrangements have also been identified in glioblastoma, AML and secretory breast cancer. In 2013, Vaishnavi et al. reported novel NTRK1 fusions in 3/91 pan-negative patients with lung adenocarcinoma using NGS and FISH (Vaishnavi et al. Nat Med. 2013 November; 19(11):1469-72)."}
{"text": "A variety of methods for preparation of polysuccinimide and subsequent hydrolysis to polyaspartic acid (or salts) have been described in the literature and patents. In addition methods of preparation of copolymers of polyaspartic acid have also been reported in the literature.\nIn a series of patents, Koskan et al. discloses a method for thermal polymerization of aspartic acid in a fluidized bed to form polysuccinimide which is then hydrolyzed to polyaspartic acid (sodium salt) using sodium hydroxide (U.S. Pat No. 5,057,597; U.S. Pat No. 5,116,513; U.S. Pat. No. 5,152,902 and U.S. Pat. No. 5,221,733). Uses of polyaspartic acid as calcium carbonate, calcium and barium sulfate and calcium phosphate scale inhibitors are also described in these patents.\nProduction of polysuccinimide and polyaspartic acid (and salts) from maleic anhydride, water and aqueous ammonia is taught in patents by Koskan and Meah (U.S. Pat. No. 5,219,952 and U.S. Pat. No. 5,296,578). Polysuccinimide is produced in at least 90% of theoretical yield by heating the maleic anhydride, water, ammonia mixture at 220.degree.-260.degree. C. In U.S. Pat. No. 4,839,461, Boemke teaches the production of a polyaspartic acid salt by the reaction of maleic acid and aqueous ammonia at 120.degree.-150.degree. C. followed by hydrolysis of the resulting acid with metal hydroxides or ammonium hydroxide. A process is disclosed (in U.S. Pat. No. 5,288,783) for the preparation of a salt of polyaspartic acid by reacting maleic acid and ammonia in a molar ratio of 1:1-2.1 at 190.degree.-350.degree. C. for a time followed by hydrolysis of the resultant polymer using metal or ammonium hydroxide.\nFox and Harada (\"A Laboratory Manual of Analytical Methods of Protein Chemistry Including Polypeptides,\" P. Alexander and H. P. Lundgren, Ed., Pergamon Press, Elmsford, N.Y., 1966, p127-151) thermally polymerized aspartic acid using 85% phosphoric acid and polyphosphoric acid to obtain improved yields and higher molecular weight polysuccinimide. Neri et al [J. Med Chem., 16, 893 (1973)] also used phosphoric acid to facilitate condensation of aspartic acid under vacuum and in a thin film process to obtain high molecular weight polysuccinimide.\nAspartic acid has been thermally copolymerized simultaneously with seventeen amino acids at 175.degree.180.degree. to obtain polymers of molecular weight range 3000-9000.[S. W. Fox and K. Harada, J. Am. Chem. Soc., 82, 3745 (1960).]Copolymers of aspartic acid with glutamic acid and terpolymers of aspartic acid, glutamic acid and alanine have also been reported in the literature (S. W. Fox and K. Harada, \"A Laboratory Manual of Analytical Methods of Protein Chemistry Including Polypeptides,\" P. Alexander and H. P. Lundgren., Ed., Pergamon Press, Elmsford, N.Y., 1966, p 127-151). Lysine has been copolymerized with aspartic acid and other amino acids and with non-amino acid monomers such as caprolactam, succinic acid, terephthalic acid (S. W. Fox and K. Harada, vida supra).\nHarada [K.Harada, J. Org. Chem., 24 1662 (1959)] reported preparation of polysuccinimide by thermally condensing precursors of aspartic acid such as monoammonium malate, monoammonium maleate, maleamic acid and combinations of asparagine and malic acid, maleamic acid and malic acid, monoammonium malate and maleamic acid, malic acid and ammonium maleamate, maleic anhydride and ammonium maleamate, fumaric acid and ammonium maleamate.\nCopolyamino acids of aspartic acid are prepared by Harada, et al (U.S. Pat. No. 4,590,260) by thermally polymerizing at least one amino acid with at least one precursor of aspartic acid such as monoammonium malate, ammonium salts of maleic or fumaric acid, an ammonium salt of maleic, malic, fumaric acid monoamide or diamide and hydrolyzing the reaction mixture under neutral or alkaline conditions.\nHarada and Shimoyama (U.S. Pat. No. 4,696,981) prepared polysuccinimide from precursors of aspartic acid such as monoammonium, diammonium, monoamide, diamide and monoamideammonium salts of malic, maleic and fumaric acid and mixtures of these materials by irradiating them with microwaves. The resulting polysuccinimide was hydrolyzed to form polyaspartic acid. Similarly mixtures of at least one amino acid and precursors of aspartic acid were irradiated with microwaves followed by hydrolysis to produce copolyamino acids of aspartic acid (Harada and Shimoyama, U.S. Pat. No. 4,696,981).\nGerman laid open document No. 4217847 discloses preparation of aspartic acid-amino acid copolymers, prepared by thermal condensation of L-aspartic acid with other amino acids followed by hydrolysis of the condensation reaction mixture.\nThese polymers were used to prevent encrustation during sugar juice evaporation.\nModified polyaspartic acids are prepared by polycocondensation of aspartic acid with carboxylic acids (monobasic and polybasic) and anhydrides, hydroxycarboxlic acids, alcohols, amines, alkoxylated alcohols and amines, amino-sugars, carbohydrates, sugar carboxylic acids, and non-protein forming amino-carboxylic acids (Ger. laid open document No. 4221875).\nU.S. Pat. No. 5,286,810 discloses the preparation of higher molecular weight copolymers of polyaspartic acid which are suitable for the inhibition of scale deposition by reacting maleic acid and ammonia in stoichiometric excess with a diamine or a triamine at 120.degree.350.degree. C. The resulting copolymers of polysuccinimide are converted to a salt of the copolymer of polyaspartic acid by hydrolysis with a hydroxide.\nCopolymers of polyaspartic acid are also made by reaction of part of the succinimide units in polysuccinimide with amines followed by hydrolysis of the remaining succinimide units to form aspartic acid units (Fujimoto et al., U.S. Pat. No. 3,846,380; Jaquet et al., U.S. Pat. No. 4,363,797).\nCopolymers of polyaspartic acid are prepared by making maleic half esters followed by addition of an equivalent of ammonia and an amine and heating to 120.degree.-350.degree. C. When an equivalent of alcohol is distilled off, a copolymer of polysuccinimide is formed which is hydrolyzed with hydroxides to form amide copolymers of polyaspartic acid (Wood, U.S. Pat. No. 5,292,858)."}
{"text": "1. Field of the Invention\nThe present invention relates to the game of pool.\nMore particularly, the present invention relates to a portable pool game.\n3. Description of the Prior Art\nEvery pool player, whether because of playing a better player or because of playing someone on a run, has unhappily experienced watching the game more than shooting in the game. In a conventional game of pool, players shot until they missed. If a player was good he could end up shooting for a long period of time.\nNow through the games of \"Super Pool\", and \"Hourglass Pool\", of the present invention, all of the old games, and the new games introduced through the addition of six new balls, can be played using equal shooting time, if desired. The option to play equal shooting time games is still present, if desired.\nNumerous innovations for the game of pool have been provided in the prior art that are adapted to be used. Even though these innovations may be suitable for the specific individual purposes to which they address, they would not be suitable for the purposes of the present invention as heretofore described."}
{"text": "Removal of an intervertebral disc is often desired if and when the disc degenerates. The disc may be replaced with a device such as a cage or other spacer that restores the height of the disc space and allows bone growth through the device to fuse the adjacent vertebrae. Spacers often do not intimately connect the two vertebral bodies and a combination of plates and screws are often used to obtain the rigidity necessary to enable bone to grow and fuse the adjacent vertebral bodies.\nImplants for spinal fusion that are impacted into the disc space and allow growth of bone from adjacent vertebral bodies through the upper and lower surfaces of the implant are known in the art. Such implants are typically provided with a lordotic taper to enable a surgeon to recreate an appropriate lordotic curvature to the motion segment. In order to create the appropriate environment for fusion, fixation hardware is applied to the spinal segment to limit the relative motion between the vertebral bodies to be fused.\nFurthermore, interbody implants that feature a screw thread form connected to a central body have been developed, such as the well-known cylindrical threaded spacers. These devices are typically hollow and allow bone growth through fenestrations in the device. Clinically, these devices are associated with the risk of post-operative loss of disc height due to the small surface area available to resist subsidence into the adjacent vertebral body relative to design of the impacted cages.\nAttempts to combine the features of the impacted implants with the implants using screw thread forms, provide a greater resistance to subsidence; however, they offer little resistance to anatomic motions where the vertebral bodies move apart from each other, such as is typical in flexion and lateral bending."}
{"text": "1. Field of the Invention\nThis invention relates to an ignition system for an internal combustion engine, and more particularly to an ignition system for an internal combustion engine of the current interruption type.\n2. Description of the Prior Art\nSuch an ignition system of the current interruption type in which a battery is used for its power supply is typically disclosed in, for example, Japanese Utility Model Publication No. 10812/1963. The conventional ignition system disclosed generally includes a transistor switch which is turned off at an ignition position of an internal combustion engine and constantly turned on at an position other than the ignition position, a pulse signal generating device for generating a pulse signal at the ignition position, and a circuit for turning off the transistor switch when the pulse signal generating device generates the pulse signal.\nIn the conventional ignition system constructed as described above, when the transistor switch is turned off, a high voltage is induced across a primary winding of the ignition coil. The so-induced high voltage is further increased by the ignition coil, so that a high voltage for ignition may be induced across a secondary winding of the ignition coil.\nIn the above-described construction of the conventional ignition system, the transistor switch which functions to control the primary current of the ignition coil is turned off for a short period of time at the ignition position, however, it is turned on at any position other than the ignition position. This causes the primary current flowing through the ignition coil to lead to the generation of much heat from the ignition coil, thereby increasing power consumption. In order to solve such a problem, an ignition system is proposed which includes a time control circuit for controlling the time for which a primary current flows from a battery to an ignition coil, as disclosed in, for example, U.S. Pat. No. 3,605,713. Unfortunately, such a conventional time control circuit is highly complicated to a degree sufficient to complicate a circuit structure of the ignition system. Thus, the conventional time control circuit is not suitable for use for an ignition system in which a decrease in the number of parts is required for a reduction of its manufacturing cost.\nAlso, the conventional ignition system is encountered with a problem that the ignition coil is overheated when the primary current continues to flow therethrough after a stop of the internal combustion engine. In view of such a problem, an ignition system is proposed which is adapted to interrupt the flow of a primary current through an ignition coil when an internal combustion engine is stopped, as disclosed in U.S. Pat. No. 3,884,208. However, the ignition system fails to control a time for which the primary current flows."}
{"text": "Numerous healthcare and cosmetic products are applied to the skin in order to provide various benefits. Such products can include, for instance, lotions, creams, moisturizers, and the like. In some circumstances, the products are intended to provide a cooling feeling or cooling sensation to the skin once applied. Existing products typically provide skin cooling by combining skin cooling agents with other substances.\nThere are several different means to impart a cooling sensation to the skin, including using evaporation, neurosensory components, or thermodynamic agents such as phase change materials. One example of a cooling agent is menthol which provides cooling in the form of a physiological or neurosensory effect on nerve endings in the human body that sense temperature. The cooling sensation from menthol is not due to latent heat of evaporation but appears to be the result of direct stimulus on the cold receptors at the nerve endings.\nThe use of phase change materials to impart cooling is discussed, for instance, in PCT International Publication No. WO 2006/007564 entitled “Cosmetic Compositions and Methods for Sensory Cooling”, which is incorporated herein by reference. In the '564 application, a skincare cosmetic composition is described in the form of a lotion that is intended for use in after-sun products, after-shave products, and body moisturizing products. The lotion is intended to create a cooling sensation on the skin by incorporating into the lotion components that absorb heat from the skin. In particular, ingredients are incorporated into the lotion that absorb heat from the skin and melt. The components have a relatively high heat of fusion which is defined in the '564 application as the heat absorbed by unit of mass of a solid chemical element at its melting point in order to convert the solid into a liquid at the same temperature. The '564 application states that the relatively high heat of fusion facilitates the absorption of heat from the skin to aid in melting the solid ingredient when applied to the skin, thereby cooling the skin temperature.\nThe use of phase change agents to impart cooling in tissues is disclosed, for instance, in PCT Patent Application No. PCT/IB2009/051515 entitled “Tissue Products having a Cooling Sensation When Contacted with Skin”. The '515 application discloses the use of a phase change agent between multiple layers of a dry tissue web with a separate hydrophobic lotion layer on the exterior surfaces of the tissue product to provide a cooling sensation. This approach is problematic since components of the hydrophobic lotion can migrate into the hydrophobic phase change agent and disrupt its ability to cool. Alternatively, the phase change agent can migrate into the lotion on the exterior of the tissue and may cause irritation to the skin.\nTherefore, a need exists for a means to effectively hold a phase change agent on or within a substrate, such as a tissue, such that it will cool the skin without allowing irritation to the skin. There also exists a need for a substrate, such as a tissue containing the composition, such that the composition can be delivered to the nose to moisturize, cool and soothe irritated noses, while holding this phase change agent within the substrate, keeping it from irritating skin."}
{"text": "Television viewers frequently search available real-time television programs by using a remote control device to change channels. The viewer could use the numeric keys on the viewer's remote control device to enter the precise number of a particular channel. After the viewer enters the digits for a particular channel, the viewer must wait until the channel change device recognizes the viewer's entry of the digits as an instruction to change channels to the particular channel. Often, the time that the viewer must wait after the time that the viewer has entered the last digit of the channel until the time that the television displays the content of the identified channel is much longer than the viewer expects.\nAlternatively, the viewer can sequentially, incrementally change channels, typically by using the up and down channel keys on the viewer's remote control device. Television viewers who incrementally change channels typically disapprove of, and become impatient with, a long lag time between the viewer's pressing an up or down channel key on the viewer's remote control device and the actual changing of the channel.\nIncremental channel change speed may be set by a particular manufacturer to occur more slowly than desired by a particular viewer. Furthermore, incremental channel change speed can depend upon the level of integration of the viewer's entertainment system components. Mixing components of different manufacturers, or even different technology advancements by a single manufacturer, in a single entertainment system can result in degradation of incremental channel change speed.\nThere is a need, therefore, for a method and apparatus that is capable of learning, for the components of a particular entertainment system, the optimal speed and timing with which the viewer's instruction to the channel change device to change channels will be recognized by the components of a particular entertainment system. Similarly, there is a further need for a method and apparatus that is capable of constructing and delivering a change channel instruction to the components of a particular entertainment system according to the optimal speed and timing with which the viewer's instruction to the channel change device to change channels will be recognized by the components of that particular entertainment system."}
{"text": "The LCD (Liquid Crystal Display) possesses many advantages of being ultra thin, power saved and radiation free. It has been widely utilized in, such as LCD TVs, mobile phones, Personal Digital Assistant (PDA), digital cameras, laptop screens or notebook screens, and dominates the flat panel display field.\nMost of the liquid crystal displays on the present market are backlight type liquid crystal displays, which comprise a liquid crystal display panel and a backlight module. The working principle of the liquid crystal display panel is that the Liquid Crystal is injected between the Thin Film Transistor Array Substrate (TFT array substrate) and the Color Filter (CF). The light of backlight module is refracted to generate images by applying driving voltages to the two substrates for controlling the rotations of the liquid crystal molecules.\nWith the constant development of the industry manufacture technology of the liquid crystal display, the cost decrease has already become one of the most development directions in the present industry. Except optimizing the manufacture process of the liquid crystal display, and developing new material to reduce the production cost, the technology of integrating the related function modules, circuits inside the liquid crystal display panel, such as directly manufacturing the gate scan drive circuit on the thin film transistor array substrate with the array manufacture process (Gate Driver on Array, GOA) to replace the external gate scan drive IC is also the popular content of the developments which many liquid crystal display panel makers compete to involve for reducing the production cost in advance.\nAs conducting the high integration liquid crystal display panel manufacture process, it requires the manufacture environment which is severer. If the particles in the manufacture environment, the metal particles generated by the production apparatus fall on the glass substrate in the manufacture process, it does not only result in the bad quality of the liquid crystal display panel but also the short circuit phenomenon in the circuit. Such short circuit phenomenon causes the image display abnormal, and the continuous usage will make the temperature of the circuit raise and even the burning of the liquid crystal display panel. If the short circuit in the circuit cannot be detected in the production, the product must have highly possible quality risk and the yield descends, which results in the customer complaints and the Call Back of the products. It leads to the massive production cost increase.\nThe temperature detection to the inside of the liquid crystal display panel according to prior art is accomplished with the temperature detection component outside the liquid crystal display panel. The result of the external temperature detection component must be restricted by the position, and cannot perform detection to the arbitrary positions inside the entire liquid crystal display panel, and the detection precision is not enough."}
{"text": "The distribution of television signals has increasingly become based on digital methods and digitally encoded forms of video and audio signals. At the same time, higher resolution (high definition TV) has become available in the market place, commensurate with larger and higher definition displays. To meet the requirement of interconnecting such high definition displays with digital signal sources such as Digital Versatile Disc (DVD) players and receivers/decoders for digital satellite and digital cable distribution of video material, a digital interface standard has evolved, known as the High-Definition Multimedia Interface (HDMI). A detailed specification for HDMI can be obtained from the “hdmi.org” website. The HDMI specification currently available and used in this application is HDMI specification version 1.3 dated Jun. 22, 2006, which is incorporated herein by reference. This HDMI standard can be employed for connecting digital video sources to digital video sinks over a cable that carries a number of digital signals and a clock signal.\nThe inherent characteristics and manufacturing imperfections of high-speed differential signaling cables such as may be used to carry HDMI signals have an adverse effect on the high-speed signals carried by the cable.\nFor example, any cable has a limited bandwidth and therefore acts as a low pass filter. The bandwidth of the cable is related to its length, the longer the cable the greater the filtering effect and the lower its bandwidth. As a result, high-frequency signals passing through the cable are attenuated, and their edges become less sharp. This leads to an increased risk of misinterpreting the received data at the receiver end of the cable, especially for long cables and high-speed data.\nFIGS. 1A-1C illustrate the effect of the limited bandwidth of a cable on the transmitted signals. FIG. 1A illustrates a high-speed signal to be transmitted through a high-speed cable, FIG. 1B shows a distorted bandwidth-limited signal received at the receiver end of the cable (before equalization), and FIG. 1C shows the received signal at the receiver end after equalization. As seen from FIG. 1B, the signal edges are slowed and short pulses are narrowed, not reaching the full transmitted amplitude.\nDifferential signaling cables are commonly used to carry high-speed digital signals in differential form, that is, pulses of opposing polarities are transmitted on the two strands of the cable. The differential signal carried over such cables may be warped, that is the two signal components (positive and negative polarities V+ and V−) are skewed in time with respect to each other (differential skew), further distorting the received signal.\nThe impact of differential skew is depicted in timing diagrams in FIGS. 2A and 2B.\nFIG. 2A shows an example timing diagram of the two single ended signal components (V+, V−) of the differential data on an HDMI channel, as it may be transmitted by an HDMI source into a cable. A timing diagram of the corresponding differential signal (Vdiff−xmit) in FIG. 2A illustrates the corresponding differential signal that is clean and easily interpreted.\nFIG. 2B shows an example timing diagram of the two single ended signal components (V+ and V−del) of the differential data on an HDMI channel, as it might be received at the end of a cable. For the sake of clarity, only the effect of the differential skew is shown in FIG. 2B. The signals V+ and V− are skewed in time with respect to each other. The negative signal component V− is delayed with respect to the signal component V+ by a differential skew delay of Td. A timing diagram of the corresponding distorted differential signal (Vdiff−rcv) in FIG. 2B illustrates that, as a consequence of the differential skew, the differential signal Vdiff−rcv is significantly distorted with clearly visible plateaus in the signal where the differential signal is zero (0). These plateau regions can only be interpreted as noise by the receiver, the result of which is to reduce the width of the window of valid data. This reduction is seen as a closure of the receive data eye and directly compromises the channel quality. The amount of differential skew delay (Td) depends on the characteristics of each individual cable, and is basically constant.\nEarlier approaches to improving cable quality so far have been limited to embedded passive equalizer circuits within the cable, which boost high frequencies of the signals attenuated in the cable. Such equalizers are fixed to compensate for a fixed cable length.\nWhile the equalization required for a given cable depends largely on the length of the cable, other characteristics of high-speed signaling cables such as the differential skew, being more random, may vary substantially between the cables.\nAccordingly, there is a need in the industry for the development of an improved high-speed signaling cable, which would provide improved signal characteristics.\nEarlier High-Definition Multimedia Interface (HDMI) signal boosters that can be used to boost HDMI signals use external power inputs, see e.g. Long Reach™ product of Gennum corporation, which has been submitted in an Information Disclosure Statement. As a result, they cannot be embedded in a standard HDMI cable. A more recent development is a stand-alone “super booster” that can be inserted inline with a cable, and is also available integrated in an HDMI cable, see references: Gefen Inc., submitted in an Information Disclosure Statement, including an advertisement of a standalone HDMI “super booster; A manual for the standalone HDMI” super booster, which has been submitted in an Information Disclosure Statement; and an advertisement for a cable with an integrated HDMI “super booster”, which has been submitted in an Information Disclosure Statement.\nThe possibility of embedding an active device within the cable is associated with a problem. Firstly, no power input may be available for such a device except through the cable, i.e. there is no provision for external power supplies. Secondly, in the case of the HDMI cable, there is not enough power available to power a simple signal regenerator, primarily because of the specification requirement to provide a termination voltage for the inputs. As a result, the embedded active device apparently cannot be powered as required.\nIn more detail, the main power requirement for an HDMI signal booster is the requirement to provide a termination voltage (3.3V) with the capability to source 12 mA for each of three HDMI inputs. The power that is available from the cable comes from a 5V line, from which a maximum current of 5 mA can be drawn (as per HDMI specification V1.3) when the sink device is active, i.e. the total available power is limited to 50 mW. The combined power requirement of the input terminations on the other hand is approximately 12 mA*3.3V*3=120 mW. Unfortunately, these requirements cannot be met in a standard HDMI cable in a simple way.\nAccordingly, there is a need in the industry for the development of an improved signal booster with an improved power control circuit for embedded cable applications based on one or more active devices, which would avoid or mitigate the above noted problem."}
{"text": "The goal of hybrid development is to combine, in a single hybrid, various desirable traits. For field crops, these traits may include resistance to diseases and insects, resistance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination, stand establishment, growth rate, maturity, and plant and ear height is important. Traditional plant breeding is an important tool in developing new and improved commercial crops."}
{"text": "In a connection-oriented communication method using a connection-type protocol typified by TCP (Transmission Control Protocol), end-to-end connection between a communication device and a correspondent communication device is first established, and then data is sequentially transmitted via the connection. The connection is released when the data transmission is completed. This type of communication method is characterized in that once a circuit (or a session) is established between communication devices; the same circuit is dedicated to the data transmission until the session is completed.\nA communication carrier, which provides a packet communication service by way of connection-oriented communication, charges one of the parties on the basis of an amount of data transmitted and received during a connection period from when the connection is established until the connection is released. For example, in a case that a communication terminal requests a content server to deliver content to the terminal, a connection is first established between the communication terminal and the content server. The communication terminal transmits a content request message via the established connection to the content server, and, in response to the request message from the communication terminal, the content server delivers content to the communication terminal via the same connection. The connection is released after the communication terminal receives the content. The carrier operating the packet communication service performs a billing process, after obtaining information required for billing, on the basis of data transmitted or received via the connection, with the data including the request message sent from the communication terminal to the content server and the content delivered to the communication terminal.\nIn such connection-oriented communication, however, there exists a problem in that communication may be discontinued before data transmission is complete. This may occur as a result of an instruction issued by one of the communicating devices or as a result of a failure of a connection link. In the case that communication is discontinued following issuance of a request message, requested content can not be fully delivered from a content server to a requesting terminal; and consequently a carrier operating the packet communication service may not be able to obtain information required for billing, as a result of which it becomes impossible for a communication charge relative to the request message to be collected."}
{"text": "1. Field of the Invention\nThe invention relates to the field of sharpening devices for chain saw cutters and particularly to a sharpening device that is mounted to the chain saw housing, and which when actuated will automatically sharpen the cutters when the cutters are in motion.\n2. Description of Prior Art\nThere have been numerous designs for saw chain sharpeners mounted directly upon the chain saw to eliminate the tedious and time consuming hand filing of the chain cutters. Such chain saws use top sharpening chains of the type disclosed in U.S. Pat. No. 3,183,973, Chain Saw With Sharpening Means, by Muir.\nA typical prior art sharpener is disclosed in U.S. Pat. No. 3,040,602, Saw Chain Sharpener And Method by R. R. Carlson. The Carlson sharpener includes a body pivotally mounted to the saw housing which contains an adjustable sharpening element mounted thereon. To sharpen the chain one need only rotate the body so that the sharpening element contacts the chain, preferably at a point where the cutters are traveling over the driving sprocket. While this design works it had several drawbacks. One of these is that the amount of force applied to the cutters is variable and thus there is always a possibility of over-sharpening and prematurely wearing the cutters. It is also subject to clogging due to debris becoming caught between the sharpening element and body.\nAnother example can be found in U.S. Pat. No. 3,495,795, Easily Dressed Sharpener by M. D. Tupper. In the Tupper design a carrier housing is movably mounted on the saw housing. A bushing is slideably mounted in the carrier. Rotatably mounted in the bushing is a shaft having a grinding stone at one end. The stone is biased away from the bushing by means of a spring. The operator pushes a knob attached to the opposite end of the shaft causing the stone to engage the cutters. The biasing means provides a resilient contact with the cutters tending to limit the force applied. Wear is compensated for by means of a separate screw which moves the carrier toward the cutters. This design is relatively expensive to manufacture, is subject to clogging, and there is no real means to totally control the force applied. Furthermore the screw has no detent means so it is difficult to know how much advancement should be made in order to compensate for wear.\nA design which eliminates the clogging problem and provides for accurate adjustment of stone feed to compensate for wear was developed by J. L. Dilworth and is disclosed in U.S. Pat. No. 4,062,253, Chain Saw Sharpener. The Dilworth design includes a stone movable from a retracted position to one in engagement with the cutters. A spring within the housing is used to bias the stone to the retracted position. An adjustment knob is provided on the opposite end which when rotated feeds the stone to compensate for wear. An internal detent is also provided which provides accurate feed control. The movement of the grinding stone is controlled by an adjustable stop. Thus having the mechanism within the housing eliminates the possibility of clogging, and the use of the detent provides accurate feed. But even here there is a possibility of over-adjusting the stone, which can result in premature wearing out of the stone and a reduction in the life of the cutters.\nOther patents of interest are U.S. Pat. Nos. 3,147,644, Sharpening Means For Chain Saws by J. W. Oreilli; 3,596,689, Saw Chain And Sharpeners by J. W. Oreilli; and 3,301,098, Chain Saw Sharpening Wheel by J. W. Oreilli.\nIt is therefore a primary object of this invention to provide a chain saw sharpener which automatically compensates for grinding stone wear.\nAnother object of this invention is to provide a chain saw sharpener that is inexpensive to manufacture.\nA further object of this invention is to provide a chain saw sharpener that does not require precise positioning in order to accomplish sharpening of the chain.\nA still further object of this invention is to provide a chain saw sharpener that is relatively insensitive to the force applied by the operator, thus reducing the possibility of prematurely wearing out the chain."}
{"text": "The present invention relates to a receiving apparatus and a display control method and, for instance, to a receiving apparatus and a display control method which enable reception of a chargeable program and display output of information relating to the chargeable program.\nIn recent years, systems have spread in which a television signal is transmitted in the form of a digital signal from satellites such as a broadcast satellite and a communication satellite and such a broadcast signal is received and viewed in each home, for instance. Since such a broadcast system can secure as many as close to 150 channels, for instance, it can broadcast many more programs than existing ground wave broadcast systems, for instance.\nIn such broadcast systems, to allow reliable selection of a desired program from among a number of programs, it is proposed to transmit an electronic program guide (EPG) which contains programs scheduled to be broadcast. The receiving side receives and displays the EPG and selects a desired program by checking it.\nFurther, such broadcast systems are scheduled to broadcast, as part of broadcast programs, chargeable programs called pay-per-view programs (hereinafter abbreviated as xe2x80x9cPPV programsxe2x80x9d). A system for providing PPV programs is one form of what is called a video-on-demand system which allows a user to immediately view a desired program or the like upon his request.\nA PPV program is transmitted from the transmission side in a scrambled state. When a user performs a procedure for purchasing a desired program by a given manipulation, the receiving apparatus side, for instance, cancels the scrambled state of a PPV program to allow viewing of the purchased program.\nGenerally speaking, when a user determines whether to purchase a PPV program, the most important factor is how long a PPV program that the user intends to view has been broadcast from its start time.\nThat is, in a state that a subject PPV program has been broadcast for only a short time or its broadcast start time has not yet occurred, a user can recognize the content of the program completely or almost completely if he purchases and starts viewing the PPV program right away. On the other hand, in a state that a subject PPV program has been broadcast for a certain time, it is difficult for the user to recognize the content of the remaining part of the PPV program if he purchases and starts viewing it at that point; it is not very meaningful for the user to purchase this PPV program.\nFor example, it is conceivable to display, by characters, the broadcast start and end times of each PPV program as time-related information of PPV programs in a purchase guide picture for purchasing of the PPV programs.\nIn this case, for example, a user determines whether to purchase a subject PPV program by recognizing how long it has proceeded in time at the present time by referring to its broadcast start and end times that are character-displayed on a purchasing picture and checking the present time. However, in this method, it is difficult for a user to recognize quickly and sensibly the degree of progress in time until the present time of a PPV program with respect to its total broadcast duration which he should recognize most appropriately. There may occur a case that a user makes an erroneous manipulation of purchasing a PPV program in spite of the fact that it has been broadcast for a certain time, with an erroneous judgment that there still remains an ample broadcast time.\nIn view of the above, it is preferable to improve the ease of operation of an interface that is used in purchasing a PPV program, that is, to improve a purchasing picture so that a user can recognize as sensibly as possible the degree of progress in time of a PPV program until the current time with respect to its total broadcast duration.\nFurther, the above-mentioned broadcast systems include systems that perform a service called xe2x80x9cnear video on demandxe2x80x9d (hereinafter abbreviated as NVOD) as well as a PPV programs providing service. For example, the NVOD is a service in which the same broadcast program is broadcast several times with delays in start time by using a plurality of channels. With this service, even if a user fails to view, from the start, a desired program that is broadcast on a certain channel, he will be able to view it from the start with a waiting time of several minutes to tens of minutes, for instance, by selecting another NVOD channel.\nTherefore, for NVOD-type PPV programs, it is preferable to construct a more effective user interface by adapting the above-mentioned purchasing picture to NVOD, that is, by adapting it so that information relating to PPV programs that are or will be broadcast on the other channels are displayed there.\nAn object of the present invention is therefore to provide a user interface that a user can use more easily in purchasing a program.\nAccording to the invention, there is provided a receiving apparatus which receives transmitted program information of a plurality of programs together with transmitted video and audio signals thereof, comprising selecting means for selecting a desired program from among the plurality of programs; and chargeable program information display control means for displaying information relating to a chargeable program that is selected by a given manipulation from among programs that are rendered selectable by the selecting means, the chargeable program information display control means presenting a graphic display indicating progress in time of the selected chargeable program at a present time.\nAccording to another aspect of the invention, there is provided a display control method for a receiving apparatus which receives transmitted program information of a plurality of programs together with transmitted video and audio signals thereof and displays program information of a selected program, the display control method comprising the steps of judging whether a chargeable program has been selected by a tuner; and presenting a graphic display indicating progress in time of the selected chargeable program at a present time if it is judged that the chargeable program has been selected.\nWith the above constitution of the invention, the progress in time of, for instance, a PPV program (chargeable program) is graphically displayed as PPV-program-related information in a PPV program purchasing picture. Therefore, for instance, an elapsed broadcast time of the PPV program at the present time with respect to its total broadcast time can be displayed so as to be recognized visually.\nWhere a selected PPV program accommodates NVOD, such information as the progress in time of each of PPV programs having the same contents as the selected PPV program and to be broadcast on different channels can also be displayed in a purchasing picture, for instance."}
{"text": "Fungal, bacterial and viral diseases in crop plants result in reduced yields and product quality and are responsible for substantial losses to farmers. For example, rice blast, an often devastating disease that occurs in most rice growing areas worldwide, is estimated to cost farmers $5 billion a year (Moffat, 1994). The disease reduces rice yield significantly, particularly in the temperate flooded and tropical upland rice ecosystems. The use of resistant cultivars is the most economical and effective method of controlling the disease. Over the last decades, much has been learned about the genetics of resistance to the blast fungus. Many major genes for resistance have been identified and widely used in breeding programs. However, the molecular mechanism of host resistance to this pathogen is mostly unknown.\nWhen a plant is attacked by a pathogen such as the rice blast fungus, it can in most cases fend off the infection by mounting a battery of defense responses (Lindsay et al., 1993). The activation of plant defense occurs upon pathogen recognition and results in the halt of pathogen ingress. Systemic acquired resistance (SAR) is one important component of this complex system that plants use to defend themselves against pathogens (Ryals et al., 1996). SAR can be triggered by a local hypersensitive response(HR) to an avirulent pathogen, which renders uninfected parts of the plant resistant to a variety of normally virulent pathogens. SAR is a particularly important aspect of plant-pathogen response because it is a pathogen inducible, systemic resistance against a broad spectrum of pathogens.\nSignificant progress has been made recently in deciphering molecular aspects of SAR. The Arabidopsis gene NPR1/NIM1 has been cloned using a map-based strategy (Cao et al., 1997; Ryals et al., 1997). Mutants with defects in NPRI/NIMl fail to respond to various SAR-inducing treatments, displaying little expression of pathogenesis-related (PR) genes and exhibiting increased susceptibility to infections. The gene encodes a novel protein containing ankyrin repeats and shows homology to the mammalian signal transduction factor IKB subclass a, suggesting that RPN1/NIM1 may interact with an NF-KB-related transcription factor to induce SAR gene expression and trigger disease resistance (Ryals et al., 1997).\nThe ankyrin repeat is a 33-amino acid motif present in a number of proteins of diverse functions including transcription factors, cell differentiation molecules, and structural proteins (Bennet, 1993). The ankyrin motif consensus sequence contains the following sequence of amino acids shown as SEQ ID NO:1:\n-D----G-TPLH-AA-------V--LL--GA-\n(LaMarco, 1991). This motif has been shown to mediate protein interactions and is usually present in tandem arrays of four to seven copies (Michaely and Bennett, 1993). Ankyrin repeat-containing proteins have been shown to have diverse functions and to be involved in protein-protein interactions. Some of these proteins in mammals are transcription-regulating proteins, such as the NF-KB, inhibitor IKB (Baldwin, A. 1996; Whiteside et al., 1997). The NF-KB/IKB signal transduction pathways are conserved in both mammals and flies. A stimulus such as IL-1 treatment or bacterial inoculation leads to activation of a signal transduction pathway because of the degradation of IKB or its homolog and the release of the NF-KB transcription factor to the nucleus to stimulate transcription (Baeuerie and Baltimore, 1996; Baldwin, 1996). In Arabidopsis, NPR1/NIM1, which is homologous to the NF-KB inhibitor IKB, controls the onset of SAR. The transcription factor targeted by NPR1/NIM1 serves as a repressor of SAR gene expression and disease resistance either by direct or indirect action (Ryals et al., 1997).\nSAR is an important plant defense mechanism against infectious pathogens. For example, evidence suggests that SAR can protect plants against rice blast disease. The SAR inducer benzo (1,2,3)thiadiazole-7-carbothioic acid S-methyl ester (xe2x80x9cBHTxe2x80x9d) was found to be effective in controlling the blast disease in field conditions.\nA gene has been isolated from blast resistant plants that encodes a novel protein containing ankyrin repeats. This gene, designated RANK1, for rice ankyrin repeats, has significant homology to the Arabidopsis gene NPR1/NIM1 and the mammalian signal transduction factor inhibitor I-KB. The RANK1 gene encodes a protein that is believed to play an important role in rice defense to the blast pathogen infection as well as to other diseases which respond through SAR. Since both the RPN1/NIM1 and RANK1 genes code for ankyrin repeats, it is believed that these repeats may be responsible for SAR induced resistance to plant disease, especially rice blast disease.\nAccordingly, the present invention provides, in one embodiment thereof, an isolated nucleic acid comprising a sequence of SEQ ID NO:2.\nIn another embodiment, the invention provides recombinant DNA expression vectors functional in a plant cell comprising a nucleic acid of SEQ ID NO:2.\nA third embodiment is a plant cell stably transformed with a nucleic acid comprising a sequence of SEQ ID NO:2.\nYet another embodiment provides a transgenic plant transformed with a nucleic acid comprising a sequence of SEQ ID NO:2.\nThe invention further provides a method of conferring resistance to disease in a monocotyledonous plant comprising stably integrating into the genome of said plant the nucleic acid having the sequence which codes for a protein comprising the ankyrin motif sequence.\nAnother embodiment of the invention provides a method of conferring resistance to rice blast disease in a monocotyledonous plant comprising stably integrating into the genome of said plant the nucleic acid having the sequence which codes for a protein comprising the ankyrin motif sequence."}
{"text": "1. Field of the Invention\nThe present invention relates to a clutch driving device, particularly a clutch driving device that operates a clutch in a horizontal motion and vertical translation method, by generating lever-sided vertical translation via a move distance of rollers horizontally moving to both sides and converting the vertical translation into operational translation for pressing a release bearing.\n2. Description of Related Art\nIn general, while manual transmissions operate, the engine power is temporarily cut before the shift gear is engaged, and then the engine power is transmitted after the shift gear is engaged, and clutches are used for this operation.\nA pedal type clutch driving device having a structure including a clutch pedal and a clutch and pressing a clutch pack connected to the clutch pedal with a force transmitted from the clutch pedal is commonly used to operate the clutch. The pedal type clutch driving device is generally used for a manual transmission.\nThe pedal type clutch driving device using a clutch pedal is a constant-close type, in which as a driver presses down the clutch pedal, the force pressing the clutch is removed, such that the gears can be shifted.\nAMTs (Auto Manual Transmission) have been developed and practically applied to vehicle, which have the convenience of automatic transmissions, in addition to the advantages of manual transmission, with the technological development.\nFurther, DCTs (Double Clutch Transmission) are called a second generation AMT by using two clutches divided from the input shaft of a manual transmission and having a clutch/gear actuator.\nThe DCT has two clutches respectively connected to an odd-numbered shift input shaft and an even-numbered shift input shaft and is classified into a wet type similar to a wet type multi-plate type and a dry type similar to a clutch, in accordance with the clutch type.\nAccordingly, the DCT is implemented in a pre-selection way that engages in advance a shift gear connected with a second clutch in traveling where a first clutch is connected with the engine, that is, shifting according to the traveling condition is performed in advance in the gear train.\nThe two clutches connected to the odd-numbered shift input shaft and the even-numbered shift input shaft are operated by a clutch driving device, which uses electric actuator for operating the clutches.\nIn general, electric actuator requires an operational structure for holding the clutch and a clutch self-opening function for ensure fail safety concept when power is cut in the vehicle.\nFIG. 5 show a clutch driving device equipped with an electric actuator having the function.\nAs shown in FIG. 5A, the clutch driving device includes a clutch 100, a release bearing 200 changing the stroke to engage clutch 100, a lever 300 lifted to pressure release bearing 200, a reciprocating body 400 lifting lever 300 while moving in the longitudinal direction of lever 300 along support plate 310, and an actuator 500 moving forward/backward reciprocating body 400 in the longitudinal direction of lever 300, using a screw rod 510.\nActuator 500 is composed of a motor.\nIn this structure, a pivot point B where reciprocating body 400 moves in the longitudinal direction of lever 300 with respect to lever 300, the total force applied to release lever 200 satisfies Ft=Fs(b/a).\nFt is the total force applied to release bearing 200, Fs is lever spring tensile force, a is the distance from the pivot point B to the point of action C, and b is the distance from the point of the force A to the pivot point B.\nThe point of the force A is the position where reciprocating body 400 is not moved by the operation of the clutch in the entire length L of lever 300, the pivot point B is the support point of reciprocating body 400 with respect to lever 300, and the point of action C is the position where lever 300 applies force to release bearing 200.\nFIG. 5B shows the relationship of force according to mechanical dynamics relationships when the clutch driving device described above operates.\nAs shown in the figure, as actuator 500 is operated to engage clutch 100, reciprocating body 400 moves forward in the length direction of the lever 300 by the rotation of screw rod 510 and pivot point B correspondingly moves.\nPivot point B moves in the longitudinal direction of lever 300 in accordance to the move distance of reciprocating body 400, such that the distance b from the point of the force A to the pivot point B increases and the total force Ft applied to release bearing 200 increases in the relationship Ft=Fs(b/a). Accordingly, the force applied to clutch 100 through bearing 200 increases and the clutch is strongly engaged.\nAs described above, a mechanical structure changing the position of the pivot point B of reciprocating body 400 in the longitudinal direction of lever 300 is used in order to variably use the lever ratio b/a in this method.\nHowever, since the change in lever ratio b/a for operating clutch 100 is made by the positional change of the pivot point B, the change in the force of action depending on the position of the pivot point B should be reflected to lever 300. Further, since the operational force of clutch 100 is achieved by forward/backward move of the pivot point B in the longitudinal direction of lever 300, energy consumption increases.\nThe information disclosed in this Background section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to photographic slide mounting apparatus.\n2. Description of the Prior Art\nPhotographic slides are produced by mounting a photographic film transparency into a slide mount frame so that the image of the photographic transparency is aligned with the aperture of the slide mount frame. A variety of different types of mounting frames and mounting apparatus have been developed.\nOne particularly advantageous type of photographic slide mount is the Pakon Slide Mount, which is a one-piece plastic slide mount sold by Pako Corporation, the assignee of the present application. The film transparency is mounted by flexing open a film insertion slot in the slide mount by means of mounting equipment. The transparency is inserted into the mount and the mount is closed. The spring-like properties of the plastic slide mount material provide the transparency with a safe and tight fit in the slide mount without the need for welding or sealing. United States Patents showing slide mounts and slide mounting apparatus of this general type include the following patents:\n______________________________________ Inventor(s) U.S. Pat. No. ______________________________________ Florjancic et al. 3,341,960 Mundt et al. 3,470,642 Mundt et al. 3,478,456 Mundt et al. 3,524,299 Mundt et al. 3,562,074 Mundt 3,570,342 Mundt et al. 3,614,854 Florjancic 3,788,031 Mundt et al. 3,807,121 Mundt et al. 3,943,029 Mundt et al. 3,977,280 Urban 4,004,340 Urban et al. 4,135,343 Thompson 4,102,029 ______________________________________\nApparatus has been developed for both manual and automatic mounting of transparencies in Pakon Slide Mounts. The manual mounting procedure utilizes a hand-held mounting device into which the slide mount is inserted. By grasping the mount and the mounter together at one side, the film insertion opening is widened to permit insertion of a transparency into the slide mount. The transparency has previously been cut from a strip or web of photographic film containing many individual transparencies and is inserted manually into the slide mount.\nWhile the hand mounting apparatus procedure is adequate for mounting small quantities of transparencies in slide mounts, it clearly is not suitable for large-scale production of mounted transparencies as is required in professional photofinishing laboratories. The Pakon 509 Slide Mounter sold by Pako Corporation is an automatic, motor-driven apparatus which mounts photographic film transparencies in Pakon Slide Mounts at rates of up to 160 slides per minute.\nIn some cases, however, the quantity of transparencies to be mounted by a photofinishing laboratory is not enough to justify the use of automatic slide mounting apparatus such as the Pakon Slide Mounter, yet is greater than that which can be efficiently performed manually. To meet this need, semi-automatic slide mounters have been developed, such as the Type 6001 and 7004 slide mounters developed by Geimuplast Peter and Mundt KG. These semi-automatic slide mounters operate in a manner generally similar to the automatic Pakon 509 Slide Mounter but are driven by an operating handle which is moved by the operator, rather than being motor driven.\nThe Type 6001 and 7004 semi-automatic slide mounters are operated by moving the operating lever through an operating cycle. During this cycle, the following five functions are performed. First, an insertion slot in a slide mount is widened to receive the transparency. Second, the film web is advanced and inserted into the mount. Third, the transparency is severed from the remainder of the film web. Fourth, the transparency is inserted completely into the slide mount. Fifth, the mounter ejects the mounted slide. These five functions form a complete mounting cycle for each transparency.\nEach transparency of a film web contains a photographic image representing a singular instant in time, which in many cases cannot be recreated if the transparency is lost or damaged. Since photographic film transparencies are such a unique commodity, it is very important that the film web and transparencies progress smoothly through the slide mounter to prevent damage to the photographic images contained thereon. If the film web becomes misaligned in the film track, it can become jammed in the machine or miscut by the knife and destroyed. This problem is amplified by the fact that photographic film is coated on one side with an emulsion and this emulsion causes the film to curl transverse to its direction of advancement through the slide mounter. In addition, photographic film strips are usually spliced together to form a \"film web\" and wound on a reel for developing so that when unwound for mounting in slide mounts, the film web has a longitudinal curl due to being wrapped around the reel. These two types of curl, transverse and longitudinal, can cause considerable problems in feeding the film web through the film track into the slide mount insertion opening in a uniform manner.\nIn addition to the problems of film curl, it is important that the film track and any areas through which the film web passes be clear of possible obstructions on which the film web could catch and be damaged. Once a transparency has been partially inserted into the slide mount and severed from the film web, it must still be fully inserted into the slide mount without damage to the transparency. The insertion opening of a slide mount must be maintained in a flexed open position as the transparency is fully inserted into the slide mount to prevent scratching or damage to the photographic image."}
{"text": "G-protein-coupled receptors (GPCRs) comprise a large super-family of integral membrane proteins characterized by having 7 hydrophobic alpha helical transmembrane (TM) domains with three intracellular and three extracellular loops (Ji, et al., J Biol Chem 273:17299-17302, 1998). In addition all GPCRs contain N-terminal extracellular and C-terminal intracellular domains. Binding of extracellular ligand may be mediated by the transmembrane domains, the N-terminus, or extracellular loops, either in alone or in combination. For example binding of biogenic amines such as epinephrine, norepinephrine, dopamine, and histamine is thought to occur primarily at the TM3 site while TM5 and TM6 provide the sites for generating an intracellular signal. Agonist binding to GPCRs results in activation of one or more intracellular heterotrimeric GTP-binding proteins (G proteins) which, in turn, transduce and amplify the signal by subsequent modulation of down-stream effector molecules (such as enzymes, ion channels and transporters). This in turn results in rapid production of second messengers (such as cAMP, cGMP, inositol phosphates, diacylglycerol, cytosolic ions).\nGPCRs mediate signal transduction across a cell membrane upon the binding of a ligand to a GPCR. The intracellular portion of the GPCR interacts with a G protein to modulate signal transduction from outside to inside a cell. A GPCR is thus coupled to a G protein. There are three polypeptide subunits in a G-protein complex: an alpha subunit—which binds and hydrolyzes GTP—and a dimeric beta-gamma subunit. In the inactive state, the G protein exists as a heterotrimer of the alpha and beta-gamma subunits. When the G protein is inactive, guanosine diphosphate (GDP) is associated with the alpha subunit of the G protein. When a GPCR is bound and activated by a ligand, the GPCR binds to the G-protein heterotrimer and decreases the affinity of the G alpha subunit for GDP. In its active state, the G subunit exchanges GDP for guanine triphosphate (GTP) and active G alpha subunit disassociates from both the GPCR and the dimeric beta-gamma subunit. The disassociated, active G alpha subunit transduces signals to effectors that are “downstream” in the G-protein signaling pathway within the cell. Eventually, the G protein's endogenous GTPase activity returns active G subunit to its inactive state, in which it is associated with GDP and the dimeric beta-gamma subunit.\nThe transduction of the signal results in the production of second messenger molecules. Once produced, the second messengers have a wide variety of effects on cellular activities. One such activity is the activation of cyclic nucleotide-gated (CNG) channels by the cyclic nucleotides cAMP and cGMP. CNG channels are membrane spanning molecules that control the flux of cations through the cellular membrane. The channels are activated-opened-by increased intracellular concentrations of cyclic nucleotide. Once opened the channels conduct mixed cation currents, including ions of Na+, K+, Mg2+ and Ca2+, for example. The activity of the CNG channels couples electrical excitation and Ca2+ signaling to changes in the intracellular concentration of cyclic nucleotides (FIG. 1).\nReceptor function is regulated by the G protein itself (GTP-bound form is required for coupling), by phosphorylation (by G-protein-coupled receptor kinases or GRKs) and by binding to inhibitory proteins known as β-arrestins (Lefkowitz, J Biol Chem, 273:18677-18680, 1998). It has long been established that many medically significant biological processes are mediated by proteins participating in signal transduction pathways that involve G proteins and/or second messengers (Lefkowitz, Nature, 351:353-354, 1991). In fact, nearly one-third of all prescription drugs are GPCR ligands (Kallal et al., Trends Pharmacol Sci, 21:175-180, 2000).\nGPCRs fall into three major classes (and multiple subclasses) based on their known (or predicted) structural and functional properties (Rana et al., Ann Rev Pharmacol Toxicol, 41:593-624,2001; Marchese et al., Trends Pharmacol Sci, 20:370-375, 1999). Most of these receptors fall into class A, including receptors for odorants, light, and biogenic amines, for chemokines and small peptides, and for several glycopeptide/glycoprotein hormones. Class B receptors bind higher molecular weight hormones while class C includes GABAB receptors, taste receptors, and Ca2+-sensing receptors. GPCRs are found in all tissues. However, expression of any individual receptor may be limited and tissue-specific. As such some GPCRs may be used as markers for specific tissue types.\nAs might be expected from the wide range of GPCRs and GPCR ligands, aberrant function of these molecules has been implicated in a large number of human disease states (Rana et al. and Ji et al., supra). GPCR agonists and antagonists have been developed to treat many of these diseases. For example the important group of receptors for biogenic amines has been the target of a large number of successful drugs. Among the receptors in this group are those for epinephrine and norepinephrine (α- and β-adrenergic receptors), dopamine, histamine, and serotonin. Examples of diseases in which GPCR function has been implicated include, but are by no means limited to: heart disease (e.g. tachycardia, congestive heart failure, etc.), asthma, hypertension, allergic reactions (including anaphylactic shock), gastrointestinal disorders, and a wide range of neurological disorders (e.g. Parkinson's disease, depression, schizophrenia, etc.). Finally, many receptors for drugs of abuse are GPCRs.\nIn many animals, GPCRs are found throughout the organism and are responsible for the maintenance of normal function as well as for pathological conditions. In other instances, the expression of specific GPCRs or families of GPCRs is very tightly controlled, e.g., being expressed only during early developmental stages, etc. Consequently, it is important to find compounds that can stimulate or activate GPCRs, or inhibit or deactivate GPCRs as needed. Agonists—compounds that stimulate the normal function of the GPCRs—have been used to treat asthma, Parkinson's disease, acute heart failure, osteoporosis, hypotension, etc. Antagonists, compounds that interfere with or block normal function have been used to treat, hypertension, myocardial infarction, ulcers, asthma, allergies, psychiatric and neurological disorders, anorexia and bulimia.\nIn addition to well-characterized receptors, many “orphan” receptors have been cloned (Marchese et al., supra) which are known from sequence similarities to be part of these families, but for which no function or ligand(s) have been discerned. Given the central role of GPCRs in control of diverse cellular activities, there remains a need in the art for methods to identify the agonists and antagonists of these “orphan” receptors as well as to identify additional antagonists for those receptors whose agonists—ligands—are known.\nAs the first recognized second messenger, cAMP is synthesized by adenylate cyclase in response to activation of many receptors coupled to G proteins Gs , and Golf and cyclase activity is inhibited by activation of receptors coupled to G protein Gi. cAMP activates cAMP-dependent protein kinase A (PKA) resulting in profound cellular responses. Physiologically, cAMP mediates such hormonal responses as mobilization of stored energy (e.g., the breakdown of carbohydrates in liver or triglycerides in fat cells, conservation of water by kidney, and Ca2+ homeostasis), control of the rate and contraction force of the heart muscle, relaxation of smooth muscle, production of sex hormones, and many other endocrine and neural processes.\nThere are a number of cAMP assays currently available. They include transcription reporter assay where a luciferase reporter is driven with a cAMP response promoter element CRE, cAMP immunoassay (Applied Biosystems Forster City, Calif.), an in vitro enzymatic assay for adenylyl cyclase (Molecular Devices, Sunnyvale, Calif.) and cAMP fluorescence polarization assay (PerkinElmer Life Sciences, Boston, Mass.). However, all these assays are end point assays where the cells are lysed and extracts are used for the tests. R. Y. Tsien and his colleagues have also developed fluorescent probes that report cAMP levels in single cells. However, the methods of application of these probes to cells makes them not suitable for high throughput screening formats (Adams et al., 1991, Nature 349:694-697; Zoccolo et al., 2000, Nat. Cell Biol. 2:25-29). There is a need in the art to be able to detect the activation of individual living cells for their cAMP production, particularly in a heterogeneous cell or tissue environment. Such detection capability would further allow the examination of receptor activation and cellular response to complex stimuli, as in the case of induced long-term memory. There also exists in the art a need for the ability to directly examine the cAMP in live cells in order to identify ligands for orphan GPCRs based on the concurrent examination of both Ca2+ and cAMP activation in a given cell as well as to identify agents that modulate GPCR-mediated activity. These and other needs are met by the present invention."}
{"text": "1. Field of the Invention\nThis invention relates to a burner for introducing a combustible mixture comprising hydrocarbon fuel, free oxygen-containing gas, and optionally a temperature moderator (liquid or vapor) into a free-flow partial oxidation synthesis gas generator.\n2. Description of the Prior Art\nIn the partial combustion of a hydrocarbon with oxygen, or air enriched with oxygen, in the presence of steam and/or carbon dioxide, temperatures between 1,100.degree. and 1,500.degree. C. are often reached. Special requirements are therefore placed on the design and the material from which the burner is constructed to avoid damage to the latter.\nAn essential requirement in burner construction of the type contemplated is that they be cooled or otherwise protected from the high temperature environment. This is most often achieved by circulating water or a similar coolant through the unit. Thus, by constructing the burner both internally and externally with coolant passages, a sufficient amount of heat transfer to the circulating cooling fluid can be achieved to minimize and stabilize the temperature which the burner itself reaches.\nNormally, the oxidizing flame which combusts the mixture, introduces the hot flame as well as the products of combustion into a reactor or generator. The latter is lined with a suitable refractory material to avoid damage as a result of the high temperatures that will be reached and sustained.\nA relatively vulnerable part of the burner is that section which is continuously exposed for extended periods of time to the high reactor temperatures. Although means have been provided for cooling internal segments of the burner, the problems which result from the high temperature still persist.\nFor example, external walls of the burner are generally surrounded with a cooling coil or the like which circulates a liquid such as water to effectuate a cooling action. Further, the lower or flame end of the burner is provided with internal passages which permit coolant to be internally circulated to maintain a desired temperature range.\nIn either instance, the forward most vulnerable face of the burner, can reach certain temperatures, or range of temperatures within which accumulations of particular slag or ash will tend to cling to the exposed burner face. Such a slag build-up will cause a reduction in burn efficiency and eventually impairment of operation and eventual unit shutdown.\nThese accumulations are prompted generally by back mixing of the combustible particles or ash as the particles enter the reactor. Here they are caught up into the violently turbulent flows of the gas as associated with the high velocity flame.\nMore specifically it is found that if the temperature on the exposed burner face is in excess of 750.degree. to 900.degree. F., ash particles will be prone to stick thereto. If, on the other hand, the temperature is kept lower than 750.degree. to 900.degree. F. on the face of the burner, the ash sticking will be substantially avoided.\nIn burners that function as required, it is found that a particle build-up along the burner face will generally commence at the lip of the discharge opening or nozzle. Thereafter, the build-up will progress radially outward from the nozzle and gradually cover a substantial segment of the exposed face. Slag will also build upon itself due to progressive insulation from cooling coil/channel.\nOne way for precluding or at least limiting this slag build-up along the burner face is to inject steam directly into the combustible mixture within the burner itself. This step will facilitate the avoidance of undesired build-ups at the discharge lip. It will not, however, completely preclude the accumulations as herein mentioned.\nFor example, the back mixing and flow of the particulated matter as a result of the turbulence immediately inside the reactor, will continue to cause or prompt a certain degree of build-up at the burner face.\nToward overcoming the above stated problems, the present invention is addressed to means for providing a fluid, dynamic shield which protects the entire burner face. The shield is provided in the form of one or more jets of a fluid such as steam, which are projected transversely of the face from a point at the burner periphery.\nA number of fluids such as steam, CO.sub.2 or even water could serve as the protective dynamic shield. For the present disclosure, however, the fluid will be considered to be steam.\nPhysically, one or more high velocity steam jets are caused to sweep the burner face. The jets first of all form a barrier which precludes the hot particles from getting to, or contacting the face. Secondly, the fluid jet is so aligned that it will flow parallel to the face, or will contact or impinge against the face preferably adjacent to the discharge lip. This creates a thermal radiation/convection shield to keep burner face below 750.degree. to 900.degree. F. Thirdly, the flow will thus clear the face of any accumulation that might be initiated.\nIt is therefore an object of the invention to provide a burner which is adapted for use in combusting a coal-oxygen mixture to achieve a partial oxidation of the gaseous product. A further object is to provide a burner of the type contemplated that is capable of withstanding undesired particulate depositions along the burner exposed face. A still further object is to provide a burner of the type contemplated wherein one or more high velocity fluid jets are directed to sweep the burner face and maintain it free of accumulated particulate matter and to concurrently protect the face by establishing a fluidized radiation/convection shield thereacross."}
{"text": "Underground junction boxes may remain buried for years. During that time they should protect their contents against the entry of ground water. It is known to completely fill underground junction boxes with a water-displacing medium such as grease. This is messy, however, both at the time the junction box is filled with grease and later if it becomes necessary to access any components or conductors inside the junction box.\nThere is a need for cost-effective, durable junction boxes suitable for use in below-grade applications. There is also a general need for through-fittings capable of sealing around a cable or the like at the point where the cable passes through a bulkhead."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the hybrid variety designated CH612435, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include greater yield, better stalks, better roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen comes to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop hybrid parent plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants and have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
{"text": "Golf is enjoyed by a wide variety of players—players of different genders, and players of dramatically different ages and skill levels. Golf is somewhat unique in the sporting world in that such diverse collections of players can play together in golf events, even in direct competition with one another (e.g., using handicapped scoring, different tee boxes, etc.), and still enjoy the golf outing or competition. These factors, together with increased golf programming on television (e.g., golf tournaments, golf news, golf history, and/or other golf programming) and the rise of well known golf superstars, at least in part, have increased golf's popularity in recent years, both in the United States and across the world. The number of individuals participating in the game and the number of golf courses have increased steadily over recent years.\nGolfers of all skill levels seek to improve their performance, lower their golf scores, and reach that next performance “level.” Manufacturers of all types of golf equipment have responded to these demands, and recent years have seen dramatic changes and improvements in golf equipment. For example, a wide range of different golf ball models now are available, with some balls designed to fly farther and straighter, provide higher or flatter trajectory, provide more spin, control, and feel (particularly around the greens), etc.\nBeing the sole instrument that sets a golf ball in motion during play, the golf club also has been the subject of much technological research and advancement in recent years. For example, the market has seen improvements in golf club heads, shafts, and grips in recent years. Additionally, other technological advancements have been made in an effort to better match the various elements of the golf club and characteristics of a golf ball to a particular user's swing features or characteristics (e.g., club fitting technology, ball launch angle measurement technology, etc.).\nDespite recent technological advances, “wood-type” golf clubs, particularly the driver, can be very difficult for some players to hit well. Accordingly, additional technological advances that improve a player's ability to get a golf ball airborne and improve the playability of wood-type golf clubs, particularly the driver, would be welcome in the golf world."}
{"text": "As the method of providing an audio signal to the user is digitized in the analog method, a wider volume area may be expressed. In addition, volume of the audio signal is diversified according to the content corresponding to the audio signal. This is because in an audio content production process, the intended loudness of the audio content may be set differently for each audio content. Accordingly, international standards groups such as the International Telecommunication Union (ITU) and the European Broadcasting Union (EBU) have issued standards for audio loudness. However, there is a problem that it is difficult to apply the standards issued by the international standards groups because the method and the standards for measuring loudness are different in each country.\nThe producers of the contents try to produce the contents with relatively loudly mixed contents and provide them to the users. This is because of the psychoacoustical characteristic which is recognized that the sound quality of the audio signal is improved when the volume of the audio signal is increased. As a result, there is a competitive landscape called a Loudness War. This causes a difference in loudness between contents or a plurality of contents internally, and the user repeatedly adjusts the volume of the apparatus on which the contents are reproduced, and it may suffer from the inconvenience. Therefore, there is a need for a technique for normalizing the loudness of the audio content for the convenience of the user using the content reproduction apparatus."}
{"text": "This invention relates to treatments for ear wax, and in particular to an improved ear wax solution.\nEar wax is produced by ceruminous glands, sebaceous glands, keratinocytes, and hair from the outer third of the human ear canal. Ear wax is composed of lipid coated epidermal cells, lipids, proteins and carbohydrates. It is very hydrophobic and not soluble in water. Ear wax functions as a protectant to the inner ear from infection, as well as a cleaning and lubricating agent for the external ear canal. However, accumulation and impaction of ear wax can cause itching, pain, hearing loss, perforated tympanum, tinnitis, dizziness, and increased risk of infection. Approximately 150,000 ear wax removals are performed weekly in America due to such otologic complications. Impaction of ear wax is the most common otologic problem encountered by physicians. It can affect up to 6% of the general population, and 20% of the geriatric population.\nExcessive ear wax is often removed in physicians\"\" offices using mechanical methods. A number of solutions have also been tried to help remove ear wax. Organic based solutions prove not to be very helpful as they appear only to soften the ear wax. Interestingly, water by itself proves to be partially effective. Alkaline solutions (e.g., solutions containing sodium bicarbonate) are somewhat more effective. However, there is still a need for an improved ear wax solution.\nThis invention relates to an improved ear wax solution. The ear wax solution includes a detergent which is effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear. The detergent is selected from anionic detergents, cationic detergents, zwitterionic detergents, ampholytic detergents, amphoteric detergents, nonionic detergents having a steroid skeleton, or mixtures thereof. The ear wax solution also includes a solvent which is water, a hydrophilic solvent, or a mixture thereof. The ear wax solution also includes an alkaline material effective to make the solution alkaline. The ear wax solution further includes an ionic additive effective to increase the ionic strength of the solution.\nIn another embodiment, the ear wax solution includes a detergent comprising a salt of a bile acid, the detergent being effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear. The ear wax solution also includes a solvent which is water, a hydrophilic solvent, or a mixture thereof. The ear wax solution also includes an alkaline material effective to make the solution alkaline.\nIn another embodiment of the invention, an ear wax formulation comprises a detergent, a polymer, and a solvent. The detergent is effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear. The detergent is selected from anionic detergents, cationic detergents, zwitterionic detergents, ampholytic detergents, amphoteric detergents, nonionic detergents having a steroid skeleton, or mixtures thereof. The polymer is effective to enhance the treatment by a mechanism involving at least one of reducing the irritancy of the detergent on the ear canal, and increasing the retention of the formulation in the ear canal and thereby reducing absorption of the formulation into the epidermal tissues of the ear canal. The solvent is water, a hydrophilic solvent, or a mixture thereof.\nIn a preferred embodiment of the invention, the ear wax solution includes one or more materials that enhance miscelle formation by the detergent. The formation of miscelles by the detergent optimizes the effectiveness of the ear wax solution.\nIn another preferred embodiment of the invention, the ear wax solution contains a plurality of detergents effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear. In some instances the use of two or more detergents together significantly improves the effectiveness of the ear wax solution.\nThe invention also relates to a method for removing ear wax from an ear. In a first step of the method, an ear wax solution is inserted into the ear in contact with the ear wax, the ear wax solution comprising: (a) a detergent effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear, the detergent being selected from the group consisting of anionic detergents, cationic detergents, zwitterionic detergents, ampholytic detergents, amphoteric detergents, nonionic detergents having a steroid skeleton, and mixtures thereof, and (b) a solvent selected from the group consisting of water, hydrophilic solvents, and mixtures thereof. In a second step of the method, the ear wax solution is maintained in contact with the ear wax for a time sufficient to treat the ear wax. Preferably, the ear wax solution is held inside the ear for at least about 30 seconds, more preferably at least about 1 or 2 minutes, to increase its effectiveness. In a final step, the ear wax solution and ear wax are removed from the ear.\nIn another embodiment of the method, an ear wax formulation is inserted into the ear in contact with the ear wax. The ear wax formulation comprises: (a) a detergent effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear, the detergent being selected from the group consisting of anionic detergents, cationic detergents, zwitterionic detergents, ampholytic detergents, amphoteric detergents, nonionic detergents having a steroid skeleton, and mixtures thereof, (b) a solvent selected from the group consisting of water, hydrophilic solvents, and mixtures thereof, and (c) a polymer effective to enhance the treatment by a mechanism involving at least one of reducing the irritancy of the detergent on the ear canal, and increasing the retention of the formulation in the ear canal and thereby reducing absorption of the formulation into the epidermal tissues of the ear canal. The ear wax formulation is maintained in contact with the ear wax for a time sufficient to treat the ear wax. Preferably, the ear wax formulation is held inside the ear for at least about 30 seconds to increase its effectiveness. Then, the ear wax formulation and the ear wax are removed from the ear.\nIn another embodiment of the method for removing ear wax, a first ear wax solution is inserted into the ear in contact with the ear wax. The first ear wax solution comprises: (a) a first detergent effective to loosen the ear wax, and (b) a solvent for the first detergent. The first ear wax solution is maintained in contact with the ear wax for a time sufficient to loosen the ear wax, then removed from the ear. A second ear wax solution is inserted into the ear in contact with the loosened ear wax. The second ear wax solution comprises: (a) a second detergent effective to remove the ear wax, the second detergent being different from the first detergent, and (b) a solvent for the second detergent. The second ear wax solution is maintained in contact with the ear wax for a time sufficient to enable the removal of the ear wax. Finally, the second ear wax solution and the ear wax are removed from the ear.\nVarious objects and advantages of this invention will become apparent to those skilled in the art from the following detailed description of the preferred embodiments.\nThe present invention relates to an improved ear wax solution that can be used at home or at the office of a physician or practitioner (e.g., an otologist) to improve ear hygiene and to ease ear wax removal.\nThe ear wax solution comprises a detergent and a solvent which is water and/or a hydrophilic solvent. The detergent is effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear. The detergent is selected from anionic detergents, cationic detergents, zwitterionic detergents, ampholytic detergents, amphoteric detergents, or nonionic detergents having a steroid skeleton. Mixtures of such detergents can also be used. The detergent can be synthetic, natural, or semi-synthetic.\nSuitable anionic detergents may include, but are not limited to, the following: long-chain (fatty) alcohol sulphates; alkali metal soaps, RCOOX, where X is sodium, potassium or ammonium, and R has a chain length between C10 and C20; alkyl aryl sulphonates; sulphonated olefins; sulphated monoglycerides; sulphated ethers; sulphated polyoxyethylated alcohols; sulphated oils; sulphosuccinates; sulphonated methyl esters; alkane sulphonates; phosphate esters; alkyl isethionates; acyl sarcosides; alkyl taurides; and fluorosurfactants. Some specific examples include sodium deoxycholate, sodium dodecyl sulphate, potassium laurate, hexadecylsulphonic acid, and sodium dioctylsulphosuccinate. In general, anionic detergents are preferred for use in the ear wax solution.\nSuitable cationic detergents may include, but are not limited to, the following: hexadecyl(cetyl)trimethylammonium, dodecylpyridinium chloride, dodecylamine hydrochloride, cetyl-trimethyl-ammonium-bromide (e.g., Cetrimide B.P.), and benzalkonium chloride.\nSuitable zwitterionic detergents may include, but are not limited to, the following: Zwittergent 3-08(n-octyl-N,N-dimethyl-3-ammonio-1-propanesulfonate), Zwittergent 3-10(n-decyl-N,N-dimethyl-3-ammonio-1-propanesulfonate), Zwittergent 3-12(n-dodecyl-N,N-dimethyl-3-ammonio-1-propanesulfonate) (Calbiochem, LaJolla, Calif.), and betaine and betaine-like detergents wherein the molecule contains both basic and acidic groups which form an inner salt giving the molecule both cationic and anionic hydrophilic groups (e.g., as disclosed in U.S. Pat. Nos. 2,082,275, 2,702,279 and 2,255,082).\nAmpholytic and amphoteric detergents can be either cationic or anionic depending on the pH of the solution. An example of an ampholytic detergent that may be suitable in the ear wax solution is N-dodecyl-N,N-dimethyl betaine. An example of an amphoteric detergent that may be suitable is alkyl dimethylamine betaine (e.g., Empigen B B from Albright and Wilson, Richmond, Va.). Other nonlimiting examples of amphoteric and ampholytic detergents that may be suitable are dodecylbeta-alanine, N-alkyltaurines, N-higher alkylaspartic acids, and the detergents sold under the trade name xe2x80x9cMiranolxe2x80x9d, and described in U.S. Pat. No. 2,528,378.\nPreferably, the detergent is selected from the category of detergents having a steroid skeleton. Anionic detergents having a steroid skeleton may include, but are not limited to, the following: sodium deoxycholate, sodium cholate, sodium taurocholate, and sodium taurodeoxycholate. Nonionic detergents having a steroid skeleton may include, but are not limited to, the following: N,N-Bis(3-D-gluconamidopropyl)cholamide (e.g., BIGCHAP, Dojindo Molecular Technologies, Gaithersburg, Md.), N,N-Bis(3-D-gluconamidopropyl)deoxycholamide (e.g., DeoxyBIGCHAP), and digitonin. Zwitterionic detergents having a steroid skeleton may include, but are not limited to, the following:, 3[(3-Cholamidopropyl)dimethylammonio]propanesulfonic acid (e.g., CHAPS). Other categories of detergents having a steroid skeleton may also be suitable.\nMore preferably, the detergent having a steroid skeleton is a natural, semi-synthetic, or synthetic bile salt. Naturally occurring bile salts are biological detergents synthesized in the liver. The commonly occurring bile acids include cholic acid, deoxycholic acid, lithocholic acid, chenodeoxycholic acid, hyodeoxycholic acid, and hyocholic acid. The bile acid can be a primary or secondary bile acid. The bile salts include alkali metal salts of such acids, such as sodium deoxycholate and sodium cholate. Most preferably, the detergent is sodium deoxycholate (xe2x80x9cDOCxe2x80x9d).\nIt is believed that bile salts such as DOC dissolve lipid bilayers of the ear wax by forming mixed micelles with lipids, and penetrate the monolayer of lipids bound to epidermal cells inside the ear. Approximately one-half of the constituents of ear wax accumulations in the ear are epidermal cells (mostly lipids and proteins); consequently, it is very important to remove these bound lipids to ease ear wax removal. It is also believed that the bile salts surround the hydrophobic parts of membrane bounded protein and move them into solution. By its dual action, the bile salt attaches itself to hydrophobic areas of the ear wax, exposing its hydrophilic tail into solution, and pulls hydrophobic particles such as membrane bounded protein into solution. The bile salt also denatures the protein.\nPreferably, the ear wax solution contains from about 0.5% to about 10% by weight of a bile salt such as DOC, more preferably from about 0.5% to about 5%. The ear wax solution acts as an ear wax softening agent at low concentrations of detergent, and an ear wax dissolution agent at higher concentrations of detergent. For example, a 1% concentration of DOC could be used to act as an ear wax softening agent, while a higher concentration of DOC could function as an ear wax dissolution agent.\nOther examples of naturally occurring detergents that may be used in the ear wax solution include phosphatides which are surface-active agents, such as lecithin and dialkylglycerylphosphorylcholine.\nIn another preferred embodiment of the invention, the ear wax solution contains a plurality of detergents effective to treat the ear wax by a mechanism involving at least one of dissolving the ear wax, softening the ear wax, and reducing the attachment of the ear wax to the ear. In some instances the use of two or more detergents together significantly improves the effectiveness of the ear wax solution. For example, the ear wax solution may contain a mixture of sodium deoxycholate and sodium dodecyl sulphate.\nIn addition to the detergent, the ear wax solution also includes a solvent for the detergent. Preferably, the solvent is water, a hydrophilic solvent, or a mixture thereof. Examples of hydrophilic solvents include alkylalcohols such as isopropanol, methanol, ethanol, n-propanol, n-butanol, secondary butanol, tertbutanol and isobutanol, alkylene glycols such as propylene glycol and polyethylene glycol, ether alcohols such as methyl cellosolve, ethyl cellosolve, propyl cellosolve, butyl cellosolve, methyl carbitol and ethyl carbitol, ether esters such as methyl cellosolve acetate and ethyl cellosolve acetate, dioxane, dimethylformamide, diacetone alcohol, methyl ethyl ketone, acetone, tetrahydrofurfuryl alcohol, and mixtures thereof. The percentage of solvent in the solution is the balance after subtracting the percentages of the other ingredients.\nPreferably, the ear wax solution also includes an alkaline material effective to make the solution alkaline. An alkaline solution increases the effectiveness of the detergent. Additionally, an alkaline solution produces an expansion of keratinocytes which provides mechanical force to disintegrate ear wax. Preferably, the solution has a pH between about 7 and about 9.5. Any suitable alkaline material can be used to make the solution alkaline. Some examples of alkaline materials include the sodium, potassium, calcium, magnesium and aluminum salts of phosphoric acid, carbonic acid, citric acid, and certain aluminum/magnesium compounds. Other examples include antacid materials such as aluminum hydroxides, calcium hydroxides, magnesium hydroxides and magnesium oxide. A preferred alkaline material for use in the ear wax solution is disodium phosphate. The alkaline material may function as a buffer in addition to increasing alkalinity. Generally, the amount of alkaline material in the solution is between about 0.1% and about 5% by weight of the solution.\nPreferably, the ear wax solution also includes an ionic additive effective to increase the ionic strength of the solution. An increased ionic strength increases the effectiveness of the detergent, for example, by insuring a large aggregation number and a small critical micellization concentration for the detergent. Any suitable ionic additive can be used in the solution. The ionic additive is preferably an alkali metal salt, more preferably an alkali metal salt of a halogen, even more preferably a chloride salt of an alkali material, and most preferably sodium chloride. Non-limiting examples of suitable ionic additives include sodium chloride, potassium chloride, sodium bromide, potassium bromide, sodium iodide, potassium iodide and the like. Generally, the amount of ionic additive in the solution is between about 0.1% and about 5% by weight of the solution. For example, in a preferred embodiment, sodium chloride is added to water in an amount to make a 0.1M sodium chloride solution.\nIn some embodiments of the invention, a polymer delivery system is added to the ear wax solution to make an improved ear wax formulation. The polymer enhances the treatment of the ear wax by reducing the irritancy of the detergent on the ear canal, and/or by increasing the retention of the formulation in the ear canal and thereby reducing absorption of the formulation into the epidermal tissues of the ear canal. These properties of the polymer allow stronger detergents to be used for ear wax removal. These properties also aid ear wax removal by increasing the ability of the formulation to wet the ear wax and the surface of the ear canal. By reducing the irritancy of the detergent on the ear canal and providing a time release delivery of the detergent, the polymer allows the use of a higher concentration of detergent without irritation of the ear. The polymer also protects and stabilizes the detergent from being broken down by any substances in the ear canal.\nA current commercial product used for ear wax removal has a tendency to develop an allergic reaction in the ear. This is not a problem with the ear wax formulation of the invention. With the addition of the polymer, the ear wax formulation can usually be left inside the ear for over 30 minutes if necessary without causing an allergic reaction. The ear wax formulation is also effective to reduce tissue inflammation and exudation in the ear.\nGenerally, the amount of polymer in the ear wax formulation is between about 0.5% and about 20% by weight of the formulation, and typically between about 0.5% and about 10%. The polymer is usually dispersed throughout the solvent. Preferably, the polymer and other materials are formulated so that the ear wax formulation has a viscosity between about 2500 cps and about 25,000 cps. Preferably, the formulation is clear to allow for better visualization of the impacted ear wax when removing it.\nThe polymer delivery system can be any polymer, or combination of polymers, capable of better retaining the formulation in the ear canal and thereby reducing absorption of the formulation into the epidermal tissues of the ear canal. The polymers can be water soluble, or non-water soluble, and can come in various lengths to accommodate one\"\"s needs. Some polymers can change from a solution state to solid state dependent upon temperature. Thus, a polymer could be in solid form at room temperature, but in a solution state when heated a few degrees more.\nPreferred polymers for use in the polymer delivery system include, but are not limited to, microsponge polymers, polytrap polymers, N,O-carboxymethyl-chitosan (xe2x80x9cNOCCxe2x80x9d), polyolprepolymers, and chitosan polymers. Microsponge polymers consist of polymeric beads having a network of pores. One such microsponge polymer is available commercially from Advanced Polymer Systems, Redwood City, Calif. The microsponge polymer is described in more detail in U.S. Pat. No. 4,690,825 to Won, issued Sep. 1, 1987, and in U.S. Pat. No. 5,145,675 to Won, issued Sep. 8, 1992 (both of which are incorporated by reference herein).\nPolytrap polymers are highly cross-linked polymethacrylate copolymers. Such a polymer is manufactured by Dow Coming Corporation, Midland, Mich., and sold under the trademark Polytrap. It is powder having particles capable of absorbing high levels of lipophilic liquids and some hydrophilic liquids. The powder structure consists of a lattice of unit particles less than one micron that are fused into agglomerates of 20 to 100 microns, and the agglomerates are loosely clustered into macro-particles or aggregates of about 200 to about 1200 micron size. Advanced Polymer Systems also sells a Polytrap System which can be used in the invention.\nNOCC is a chitosan derivative having carboxymethyl substituents on some of both the amino and primary hydroxyl sites of the glucosamine units of the chitosan structure. One such polymer is available commercially from Chitogenics, Inc., Halifax, Nova Scotia, Canada. This polymer is described in detail in U.S. Pat. No. 4,619,995 to Hayes, issued Oct. 28, 1986, U.S. Pat. No. 5,679,658 to Elson, issued Oct. 21, 1997, and U.S. Pat. No. 5,888,988 to Elson, issued Mar. 30, 1999 (all of which are incorporated by reference herein).\nPolyolprepolymer is a mixture of liquid hydroxyl terminated polymers and polyethylene glycol. One such polymer, Polyolprepolymer-2, is available commercially from Barnet Products, Inc., Englewood Cliffs, N.J.\nChitosan is deacetylated chitin, or poly-N-acetyl-D-glucosamine. It is available commercially from many sources, such as Protan Laboratories Inc., Redmond, Wash. As used herein, xe2x80x9cchitosanxe2x80x9d includes chitosan, inorganic or organic salts of chitosan, and any chemically modified forms of chitosan or chitosan derivatives. This polymer is described in detail in U.S. Pat. No. 5,141,964 to Noel, issued Aug. 25, 1992, and U.S. Pat. No. 5,744,166 to Illum, issued Apr. 28, 1998 (both of which are incorporated by reference herein). Other types of polymers can also be used in the invention.\nIn a preferred embodiment of the invention, the ear wax solution includes one or more materials that enhance miscelle formation by the detergent. The formation of miscelles by the detergent optimizes the effectiveness of the ear wax solution. Any suitable material(s) can be used to enhance miscelle formation. In some embodiments, the materials are alkaline materials such as alkaline buffers, ionic additives and/or polymers as described above. The ear wax solution may also include one or more additives (e.g., polymers or alcohols) to increase patient comfort.\nThe ear wax solution, or the ear wax formulation with the polymer delivery system, can also include a topical therapeutic agent for the treatment of the ear. Some nonlimiting examples of therapeutic agents are anti-infectives, antiinflammatory agents, analgesics, and anesthetics. When the therapeutic agent is used in an ear wax formulation containing a polymer, the polymer typically enhances the formulation by either reducing side effects of the therapeutic agent, increasing the therapeutic efficacy of the therapeutic agent, or improving the stability of the formulation.\nThe ear wax solution of the invention can be prepared in any suitable manner. Typically, the solution is prepared by adding the ionic additive to the solvent in a desired concentration and mixing the solution, and then adding the desired amounts of alkaline material and detergent and further mixing the solution until the materials are dissolved. The ear wax formulation containing the polymer can also be prepared in any suitable manner. Typically, the formulation is prepared by initially mixing the detergent with the polymer so that it is incorporated into the polymer, and then adding the detergent/polymer to a solution prepared as described above.\nIn another embodiment of the method for removing ear wax, the ear is washed with a first ear wax solution to loosen the ear wax, and then a second ear wax solution is used to remove the ear wax. In this method, the first ear wax solution is inserted into the ear in contact with the ear wax. The first ear wax solution comprises: (a) a first detergent effective to loosen the ear wax, and (b) a solvent for the first detergent. The first ear wax solution is maintained in contact with the ear wax for a time sufficient to loosen the ear wax, then removed from the ear. The second ear wax solution is inserted into the ear in contact with the loosened ear wax. The second ear wax solution comprises: (a) a second detergent effective to remove the ear wax, the second detergent being different from the first detergent, and (b) a solvent for the second detergent. The second ear wax solution is maintained in contact with the ear wax for a time sufficient to enable the removal of the ear wax. Finally, the second ear wax solution and the ear wax are removed from the ear."}
{"text": "The present invention relates to a multichip semiconductor device using multiple chips.\nThe present invention also relates to a chip for a multichip semiconductor device and a method of manufacture thereof.\nRecent computers and communication equipment use for their important section a large-scale integrated circuit (chip) which has a great number of electronic components, such as transistors, resistors, etc., integrated into a semiconductor substrate. Thus, the performance of the entire equipment depends largely on the performance of the chip.\nOn the other hand, so-called multichip semiconductor devices have proposed, each having a plurality of chips to improve the whole performance of the equipment. FIGS. 1, 2 and 3 are sectional views of conventional multichip semiconductor devices.\nFIG. 1 shows a multichip semiconductor device of a type in which a plurality of chips 82 are placed side by side on a multilayered interconnection substrate 81. Reference numeral 83 denotes a solder bump.\nFIG. 2 shows a multichip semiconductor device of a type in which chips are connected together with their major surfaces opposed to each other. FIG. 3 shows a multichip semiconductor device of a type in which a plurality of chips 82 are stacked using stacking plates 84.\nHowever, these conventional multichip semiconductor devices have the following problems.\nIn the multichip semiconductor device shown in FIG. 1, the plane area of the device increases because the chips 82 are arranged in the same plane.\nThe conventional semiconductor device of FIG. 2 is free of the problem with the device of FIG. 1 that the plane area of the device increases. This is because the chips 82 are stacked one above another. However, the device of FIG. 2 has a problem that the number of chips that can be stacked is limited to two. In addition, it is difficult to electrically test each chip.\nThe conventional semiconductor device of FIG. 3 does not suffer from the problems with the conventional semiconductor devices of FIGS. 1 and 2. However, its structure is complex, its thickness is great, and its manufacturing cost is high. This is because a stacking plate 84 need to be provided between any two adjacent two chip."}
{"text": "Environmental contamination with radioactive materials is a common problem. The problem may occur as a result of mining operations, such as for uranium, or contamination due to operation of nuclear facilities with inadequate environmental controls, or from the disposal of radioactive wastes. Alternatively, contamination may occur as a result of dispersion of uranium billets which have been used as a high density material in military or civil applications as a result of warfare or civil accident.\nMining operations have established practical and economic methods for the economic recovery of some radioactive elements from contaminated materials. The objective of mining, however, is usually the economic recovery of materials and secondary waste is rarely the major issue. In environmental clean-up, the economic objective is to complete effective clean-up with minimum secondary waste at minimum cost, and the value of recovered radioactive substances is of secondary importance. Techniques and chemicals which would not be economical or appropriate for mining applications may become practical for environmental clean-up.\nIt is well established that radioactive elements can be recovered from environmental materials by mechanically washing with water with or without surface active additives. However, such procedures are generally limited to the mechanical separation of solids, and will not remove contaminants that are chemically bound to the solid phase.\nThere are established chemical methods for dissolving insoluble radioactive contaminants in concentrated solvents, such as strong acids, in a process known as acid leaching. Such procedures are effective, but are disadvantageous if the spent concentrated solution ultimately becomes waste. In many cases, the concentrated solvents themselves are hazardous in addition to containing the radioactive contaminant that the process is designed to concentrate. The acid leaching and other processes using concentrated solvents to dissolve the radioactive contaminant have the further disadvantage of also dissolving other contaminants that the process was not designed to remove, such as nonradioactive metals.\nIn the decontamination of internal surfaces of nuclear reactor circuits, early processes involved washing with concentrated chemical solutions to dissolve contaminants to yield a concentrated solution containing the contamination. The processing of these waste solutions was found to be difficult and inconvenient and resulted in them becoming waste and requiring disposal. The technology has now progressed to allow the recovery of radioactivity, typically by ion exchange, in a dilute acidic recirculating system. These solutions, being dilute and acidic, do not contain carbonate and are not particularly useful or appropriate for dissolving actinide elements because they do not form soluble complexes with the actinide elements.\nIn reactor decontamination processes, it has been established that certain organic reagents can be used to dissolve contamination and yield it to an ion exchange resin in a recirculating process in such a way that the organic reagent is continuously re-used. Examples of solutions used in acidic reactor decontamination processes are vanadous formate, picolinic acid and sodium hydroxide. Other processes typically use mixtures of citric acid and\noxalic acid. These reactor decontaminating solutions have the disadvantage of not being capable of being used in a single one time application to dissolve actinides, radium, and certain fission products, such as technetium.\nPrevious reactor decontaminating solutions do not contain carbonate and are acidic, dissolving the iron oxides of the radioactive elements commonly found in contaminated reactor circuits. This nonselective metal dissolving capacity is a disadvantage of the acidic solutions and makes them unsuitable for decontamination of material such as soil that contains iron and other metals that are not intended to be recovered. Another disadvantage of acidic solutions is that materials such as concrete or limestone are subject to damage or dissolution in an acidic medium. Also, in dealing with previously known washing solutions for treating soil, these solutions contain too many nonselectively dissolved contaminants preventing subjection of the solution to recovery of contaminants and recirculation of the solution to accomplish further decontamination.\nIt has been established that uranium and transuranic radioactive elements can be dissolved in concentrated acidic (pH&lt;1) chemical systems. The acidity poses difficulties as discussed above. Uranium and sometimes thorium are recovered in mining operations in a concentrated basic medium containing carbonate. The use of concentrated solutions is motivated by the need to dissolve materials at a rate economic for mining operations, and such solutions are not particularly suitable where avoidance of secondary waste is of primary concern. There are also references that suggest that uranium and plutonium can be dissolved in a dilute basic solution containing carbonate, citrate (as a chelating agent) and an oxidizing or reducing agent. Such solutions are not, however, suitable for the recovery of radium/barium sulfate because they do not form soluble complexes from barium sulfate."}
{"text": "1. Field of the Invention\nThis invention relates to an image holding member for electrophotography sensitive to electromagnetic waves such as light (here referred to a light in a broad sense including ultraviolet ray, visible light, infrared ray, X-ray, .gamma.-ray and the like).\n2. Description of the Prior Art\nHeretofore, as photoconductive materials constructing image holding members for electrophotography, there have been known generally inorganic photoconductive materials such as Se, Cds, ZnO and the like and organic photoconductive materials (hereinafter referred to as \"OPC\") such as PVK (polyvinylcarbazole), TNF (tetranitrofluorenone) and the like. However, at present the known photoconductive materials can not simultaneously satisfy the characteristics required for image holding members for electrophotography such as high sensitivity, good S/N ratio (photo-current/dark current), wide range absorption spectrum, high light response, high dark resistance, stability upon repeated use and the like.\nThere have been recently proposed various practical OPC type image holding members since they are of low cost and do not cause any environmental pollution, and as a result of improvement, they have now excellent electrophotographic fundamental characteristics such as latent image formation, photosensitivity and the like which there is much room for improvement as to stability upon repeated use such as resistance to mechanical, optical or electrical external action. In particular, with respect to mechanical action, the surface of the OPC type image holding member is easily scratched even with a slight action force such as, for example, the action force of a blade used for cleaning the remaining toner (such cleaning is essential to electrophotographic processes) since the surface of the image holding member is made of an organic material. While the member is repeatedly used, the toner attaches to the scratched portions and therefore, the copied images become unclear gradually. Thus, the image holding member can be used only a limited number of times.\nIn order to supplement the mechanical strength of the OPC image holding member, there is proposed to form a surface protective layer of various kinds, for example, laminating with an organic material layer having a higher hardness by coating, depositing an inorganic layer by means of vacuum vapor deposition, laminating with a material capable of being surface-hardened by applying a thermal hysteresis, and the like.\nHowever, these methods can give a sufficient function of a protective layer, but there are caused side effects undesirable for fundamental characteristics of electrophotography such as lowering of resolution, lowering of image quality, change of characteristics during using for a long time and the like, and when it is tried to make the characteristics compatible, all the characteristics often become insufficient resulting in impractical image holding members.\nAs mentioned above, any satisfactory results are not obtained. This results from the fact that there are very few materials which are excellent in mechanical strength and satisfy the function of the most surface layer of electrophotographic image holding members with respect to electric and optical characteristics. Even if there are obtained the above-mentioned desirable materials, there is not easily available a means suitable for laminating the material to form freely a thin and uniform layer without damaging the photoconductive layer. This is one of the reasons retarding solving the above problems.\nAmong the various methods proposed heretofore, the most popular method is a coating method in which a constant viscosity of a coating slurry is required so as to coat uniformly, but when a high viscosity is employed, the resulting coated layer is apt to be thick. Most of the materials having high hardness are non-photoconductive so that they retard photoconductive action and the image quality is lowered when the protective layer is unnecessarily thick.\nAlternatively, for example, when an inorganic thin film such as SiO.sub.2 and the like is produced by means of a vacuum vapor deposition, the substrate temperature should be elevated to some extent to produce a good quality film, but OPC layers are usually weak to heat and are subjected to thermal damage upon forming a protective layer. In addition, vacuum vapor deposition, sputtering and the like have a directivity with respect to film forming atoms incident upon the substrate surface so that uniform and large area can be disadvantageously obtained with difficulty."}
{"text": "Aerogels are a fascinating class of high surface-area, mechanically-robust materials with a broad range of both commercial and fundamental scientific applications. Owing to its highly porous mass-fractal nanostructure, amorphous silica aerogel has been used as a capture agent in NASA's cometary-dust retrieval missions, to control disorder in 3He-superfluid phase transitions, in the fabrication of targets for laser inertial confinement fusion, in low-k microelectromechanical (MEMS) devices, and in Cherenkov nucleonic particle detectors.\nIn particular, amorphous carbon aerogel has received a considerable amount of attention in recent years owing to its low cost, electrical conductivity, mechanical strength, and thermal stability. Numerous applications have been explored for this material including water desalination, electrochemical supercapacitors, and thermal insulation.\nImpressive advances have been made in the synthesis of polycrystalline aerogels through the oxidative aggregation of chalcogenide quantum dots that preserve spectral signatures of quantum confinement. Also, silicon divacancies in nanodiamond have also been shown to be bright single-photon-emitters at room temperature (Jelezko, Phys. Stat. Sol. A, 2006), as well as being photostable near-infrared biocompatible fluorophores (Lu, PNAS, 2007).\nFurthermore, recent high-pressure, high temperature (HPHT) experiments with mesoporous silica have been employed to produce mesoporous coesite phase after oxidative removal of the carbon pressure medium. However, the achievement of an amorphous to crystalline phase transition in an aerogel material has remained an outstanding challenge, largely due to the difficulty in preventing pore collapse in the high surface area aerogel starting material.\nIn addition, thermal, electrical, optical, mechanical, and chemical properties of low-density amorphous aerogels can change profoundly through conversion from an amorphous to a crystalline phase, opening up new horizons for applications of this material in fundamental science.\nTherefore, a method and system capable of synthesizing crystalline aerogel materials from amorphous aerogel precursors would have great utility in basic science and commercial applications."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for preparing a therapeutic agent from placenta.\n2. Description of the Prior Art\nTo date, many publications have been concerned with methods of extracting physiologically or pharmacologically active substances or purifying the placental extract. Some of these extracts have been reported to possess anti-ulcer or anti-cancer effects. None, of the reported prior art placental extract preparations have been shown to possess leukemia therapeutic effects.\nHieda (\"REIZOTAIBAN NO SEIKAGAKU TO IRYOKOKA\" or \"The bio-chemistry and therapeutic effect of frozen placenta\" Kinbara Shoten, 1965) reported the presence of a substance in human placenta which is effective against cirrhosis. Kumura (HIROSHIMA IGAKU 22 (12), 1136, 1969) and Saito (Clinical report 3 (7), 543, 1969) described a placenta preparation which possesses anti-ulcer effects. Also, Byong Ho Chin (Abstract of Papers of the 9th International Cancer Congress, 467 pp, 1966) communicated an anti-Ehrlich sarcoma agent and an anti-N-F sarcoma agent from human placenta.\nHieda's preparation against the cirrhosis was obtained as follows:\nA fresh placenta was washed with water, permitted to stand at 2.degree. - 4.degree. C for several days, minced, and boiled for 60 min. The preparation was admixed with 1 N HCl in an amount of 1/5 its volume, to obtain a pH of 1.8. It was then digested with 2 g of pepsin at 38.degree. C for 20 hours. The digested fluid was centrifuged at 3,000 r.p.m. for 15 min. to separate the supernatant from the precipitate. The supernatant was passed through an ion exchange resin to reduce its acidity to pH 4.4 - 4.6 and its volume was increased to 100 ml for every 100 g (wet weight) of placenta by the addition of water. This preparation was referred to as Solution A. The precipitate was hydrolysed with concentrated hydrochloric acid by heat treatment over a period of 10 hours. The hydrolysate was succeedingly decolorized with activated carbon, the excess volatile acid removed by evaporation on a water bath, and then subjected to secondary decolorization. The acidity of the solution was reduced with an ion exchange resin to pH 4.4 - 4.6. The volume of the eluate was made up so that every 100 g (wet weight) of placenta gave the volume of 25 ml. This preparation was referred to as Solution B. Solutions A and B were blended and the acidity of the mixture was adjusted to pH 6.1 - 6.4. Following boiling and clarification by filtration, the preparation was poured into an injection ampoule, and the solution was sterilized for use. The preparation was a transparent and yellow-colored solution, had a specific gravity of 1.0090 - 1.0132, pH of 6.1 - 6.4, showed a negative sulfosalicylic acid test and contained 78.6 - 82.3 mg/ml in dry matter. The ash, total nitrogen and amino nitrogen contents were 8.0 - 9.3 mg/ml, 9.13 - 11.21 mg/ml and 8.56 - 10.24 mg/ml respectively. The extract was claimed to possess lipotropic activity and to be capable of enhancing tissue respiration of the liver, stimulation of the thyroid gland, and basal metabolism of castrate animals. It was also reportedly useable for cirrohosis therapy, in humans and for experimental standard test animals.\nBoyous Ho Chin prepared an emulsion of placenta and centrifuged it to obtain the supernatant. With addition of alcohol, a precipitation was obtained, which was subjected to fractionation by means of paper electrophoresis. The resultant fraction was dialysed against water and the dialysate or non-dialysable fraction was further subjected to precipitation with acetone. The author claimed that this precipitate possessed inhibitory effects on the growth of Ehrlich sarcoma and N-F sarcoma.\nIt is of consequence to note that none of the previously prepared placental extracts were found to possess therapeutic effects against leukemia, as shown by the description vide infra.\nOn the other hand, Carbazilquinone (Arakawa, M. et al Gann, 61 485, 1970), cytosine arabinoside (Talley, K. et al. Blood 21 352, 1963), Daunomycin (Tan, C. et al, cancer 20 333, 1967), Adriamycin (Di Marco, A. et al. Cancer Chemotherapy reports 53 33, 1969), L-asparaginase (Kidd, G. G. et al. Journal of Experimental Medicine 98 565, 1953) are among known anti-leukemic agents. These materials are extracted from natural sources, however, other than the placenta, or they have been chemically synthesized. However, all the anti-leukemic agents, so far known and used, are not specific in effect to the leukemic cells, but are observed clinically to give rise to various undesirable side effects, such as leucopenia, thrombocytopenia, anemia, hemorrhage, vomitting, diarrhea, fever, renal lesion, hepatic lesion, jaundice, etc. Therefore, therapeutic use of these prior art materials also required auxiliary care in order to prevent such inevitable complications. Acceptable therapeutic results have, therefore, not been fully achieved, due to the quite serious complications derived from the side effects. Chemotherapy for leukemia at the present time is therefore, only capable of leading to a remission of the disease, but this is achieved by a trade-off between longer survival and host toxicity, but complete cure has not been possible."}
{"text": "Conventionally, receiving devices, which receive radio waves including desired waves and interfering waves via antennas, carry out analog-to-digital conversion, extract desired waves by use of received signal strength indicators (RSSI), carry out digital demodulation, and carry out automatic gain control (AGC) to adjust the input levels of demodulators, have been developed and applied to base-station receiving devices in cellular and mobile communication systems, base-station receiving devices in fixed-line telecommunication networks, grand-station receiving devices in satellite communication systems, and broadcasting systems.\nPatent Literature Document 1 discloses an automatic gain control method in a mobile orthogonal frequency division multiple access (OFDMA) network, which is characterized in that the power level of an analog baseband signal is adjusted based on the average power of cyclic prefixes which are calculated based on the received signal intensity. Patent Literature Document 2 discloses a receiving device which is able to substantially expand the dynamic range with the input and the output of an RSSI circuit, receive a wide range of signal levels, and transmit accurate reception levels. Patent Literature Document 3 discloses a receiving device preventing intermodulation, which reliably suppress intermodulation waves included in reception channels so as to effectively prevent a reduction of reception quality due to intermodulation. Patent Literature Document 4 discloses a broadcasting receiving device which effectively attenuates signal levels of interfering stations by use of AGC and RSSI so as to suppress received desired frequency signals from being reduced in levels. Patent Literature Document 5 discloses a receiving device which determines interfering waves based on RSSI of desired waves so as to change the linearity of a low-noise amplifier (LNA), thus reducing influences of interfering waves. Patent Literature Document 6 discloses an AGC circuit which applies an optimum gain to the power of desired waves irrespective of interfering waves being input to a receiver. Patent Literature Document 7 discloses a wireless cellular communication terminal which prevents wireless characteristics from being degraded by changing the capacitance of a capacitor connected to a detection and rectifier circuit depending on the existence or nonexistence of interfering waves or the magnitude of interfering waves. Patent Literature Document 8 discloses a receiver which carries out AGC and digital demodulation processes. Patent Literature Document 9 discloses a receiving device which carries out a software demodulation process without using AGC irrespective of fluctuations in reception levels. Patent Literature Document 10 discloses a receiver having an AGC circuit adapted to a digital CATV tuner, which is characterized by gaining optimum AGC characteristics irrespective of deviations of tuner's gains or deviations of gains between channels. Patent Literature Document 11 discloses a wireless receiving device which is able to reliably amplify desired signals while suppressing increasing power consumption at an RF frontend, which is characterized by using a DC-DC converter which controls power consumption of a low noise amplifier (LNA) based on the received signal strengths before and after a low pass filter. Patent Literature Document 12 discloses a receiving device employing a direct conversion system or a Low-IF system. Patent Literature Document 13 discloses a receiving module of a terrestrial digital television receiver which is able to suppress interfering waves due to unwanted waves of adjacent channels by applying AGC to an RF amplifier circuit based on the signal level after IF tuning. Patent Literature Document 14 discloses a receiving device having a high resistance against interfering waves, which determines the existence or nonexistence of interfering waves based on differences between RF input levels and IF input levels, thus improving a reception efficiency.\nFIG. 4 shows a block diagram of a conventionally-known receiving device (see Patent Literature Document 2). The receiving device includes an antenna 1, a mixer 3, a local oscillator 4, an A/D converter 5, an analog variable gain function part 6, a channel selecting filter 8, a received signal strength indicator (RSSI) 9, a comparison controller 11, a demodulator 12, and a low noise amplifier (LNA) 13 with a variable gain function. The receiving device prevents saturation and secures backoff (i.e. a process needed to transmit modulation signals having peak components without distortions) in the analog-to-digital (A/D) converter 5 based on the direct conversion system, the single conversion system, or the multiple conversion system.\nFor this reason, the comparison controller 11 determines the level detected with the RSSI 9 (i.e. a circuit used to monitor the desired-wave level at the output side of the channel selecting filter 8) which is arranged at the output side of the A/D converter (ADC) 5, wherein the analog variable gain function part 6, arranged at the input part of the A/D converter 5, controls a gain reduction on desired waves exceeding the allowable threshold.\nFIG. 5 is used to explain problems of the conventionally-known receiving device, wherein FIG. 5(a) is used to explain the level diagram and the AGC behavior in the case of the receiving device solely inputting desired waves, while FIG. 5(b) is used to explain the level diagram and the AGC behavior in the case of the receiving device inputting both low desired waves (close to the minimum sensitivity) and strong interfering waves.\nWhen the level of a desired wave is increased from the minimum sensitivity in FIG. 5(a), the comparison controller 11 determines the level detected with the RSSI 9 in the AGC loop based on the level of a desired wave by use of a threshold, and therefore the analog variable gain function part 6 controls a gain reduction (−a [dB], −b [dB]) on a desired wave exceeding the threshold level in order to aim to prevent saturation and secure backoff in the A/D converter 5.\nSince the level of a desired wave input to the demodulator 12 becomes higher than the minimum definition level in FIG. 5(a), it is possible to demodulate the desired wave while maintaining the carrier-to-noise (C/N) ratio to the noise level, which is needed to demodulate the desired wave. When a desired wave is solely input to the receiving device, it is possible to reduce the dynamic range which is required as the input level of the demodulator 12.\nHowever, the following problem is raised when the receiving device receives a strong interfering wave and a low desired wave (close to the minimum sensitivity) as shown in FIG. 5(b). That is, when an interfering wave instead of a desired wave becomes dominant due to the increasing level of an interfering wave, the comparison controller 11 determines the level of an interfering wave detected with the RSSI 9 in the AGC loop, and therefore the analog variable gain function part 6 controls a gain reduction (−a [dB], −b [dB]) on the interfering wave whose level exceeds the threshold in order to aim to prevent saturation due to the interfering wave and to secure backoff in the A/D converter 5. In this case, the analog variable gain function part 6 controls a gain reduction (−a [dB], −b [dB]) on a low desired wave close to the minimum sensitivity similar to an interfering wave, and therefore the level of a desired wave will be decreased in the circuit portion subsequent to the A/D converter 5.\nAdditionally, due to a desired wave and an interfering wave passing through the channel selecting filter 8 following the analog-to-digital (A/D) converter 5, the level of an interfering wave is extremely attenuated and suppressed to be lower than the level of a desired wave, wherein the level of a desired wave input to the demodulator 12 remains lower than the minimum definition level as shown in FIG. 5(b), and therefore the original information is lost to cause extreme degradation of demodulation or disability of demodulation because the information of a desired wave below the minimum definition level of the demodulator 12 is cut out. the channel selecting filter 8 following the A/D converter 5, the level of an interfering wave is extremely attenuated and suppressed to be lower than the level of a desired wave, wherein the level of a desired wave input to the demodulator 12 remains lower than the minimum definition level as shown in FIG. 5(b), and therefore the original information is lost to cause extreme degradation of demodulation or disability of demodulation because the information of a desired wave below the minimum definition level of the demodulator 12 is cut out.\nWhen the conventionally-known receiving device is used to receive and demodulate a low desired wave close to the minimum sensitivity along with a strong interfering wave, it is necessary to extremely increase the dynamic range necessary for the input level of the demodulator while increasing the number of bits necessary for the demodulator, thus increasing the circuit scale of the receiving device while increasing power consumption.\nFIG. 6 shows a block diagram of another conventionally-known receiving device (see Patent Literature Document 1). The receiving device includes an analog block 22, an antenna 24, and a digital baseband part 32. The analog block 22 includes a band-pass filter (BPF) 26, a low noise amplifier (LNA) 28, a local oscillator (LO) 30, and amplifiers (VGA) 34, 36. The digital baseband part 32 includes a received signal strength indicator (RSSI) 38, a control logic part 40, and A/D converters 46, 48. Reference sign 42 shows enable pulses.\nPatent Literature Document 4 discloses the receiving device which monitors and compares the levels before and after a channel selecting filter, which is arranged to suppress interfering waves, so as to compare the level difference and the allowable threshold, thus determining the existence or nonexistence of an interfering wave. Specifically, it is possible to determine the existence of an interfering wave due to a large difference between the levels before and after the channel selecting filter but to determine the nonexistence of an interfering wave due to a small difference between the levels.\nFor example, Patent Literature Document 4 discloses the broadcasting receiving device which determines that a small amount of radio waves (i.e. interfering waves) are being received via channels other than the currently serving channel due to a small difference between the levels before and after IF-BPF, and therefore it performs AGC via the current channel by controlling the gain of the variable-gain RF-AMP, preceding MIXER, based on the RSSI level before IF-BPF. On the other hand, the broadcasting receiving device determines that a large amount of radio waves are being received via channels other than the currently serving channel due to a large difference between the levels before and after IF-BPF, and therefore it performs AGC collectively via the other channels and the current channel by controlling the gain of the variable-gain RF-AMP, preceding MIXER, based on the level representing a plurality of reception channels which is calculated by subtracting the RSSI level after IF-BPF from the RSSI level before IF-BPF. Herein, it is possible to achieve an AGC function by applying a variable gain to RF-AMP in connection with desired waves which are detected with respect to the current channel alone or a plurality of channels including the current channel.\nThe present invention, which will be discussed later, detects the existence of a strong interfering wave based on a level difference of RSSI before the channel selecting filter while reducing the gain of the analog variable gain function part (i.e. a first AGC-controlled subject) preceding the A/D converter such that the output of the A/D converter, serving as an AGC-controlled subject in the minimum sensitivity reception and preceding channel selection, will not be distorted due to a strong interfering wave. In this case, the present invention increases the gain of the digital variable gain function part, serving as a second AGC-controlled subject, so as to correct the low level of a desired wave with the demodulation-enabled level.\nPatent Literature Document 8 imports an idea into the AGC circuit including RSSI to receive desired waves with a wide range of levels even when high-level interfering waves are intermixed with desired waves. That is, it is possible to appropriately control the AGC gain entirely over the circuitry by activating an AGC operation by use of the processing result of the broadband RSSI1 based on the output of the A/D converter (ADC) reflecting desired waves and interfering waves and the processing result of the narrowband RSSI2 solely reflecting desired waves after being subjected to the band limitation of BPF following ADC. In this connection, Patent Literature Document 9 and Patent Literature Document 10 disclose prior art which aim at improvement of Patent Literature Document 8.\nPatent Literature Document 8 discloses a method of extracting data by use of ADC at two points in different gain stages arranged in the reception IF system depending on the input level in consideration of the necessity of two systems for ADC.\nPatent Literature Document 10 discloses a method of independently varying and controlling the AGC gain of RF (radio frequency) and the AGC gain of IF (intermediate frequency) in consideration of the gain control reference solely reflecting desired waves precluding interfering waves.\nThe first embodiment of Patent Literature Document 8 refers to the circuitry in which the broadband RSSI1 based on the ADC output reflecting desired waves and interfering waves is not interlocked with the narrowband RSSI2 based on the BPF output after ADC solely reflecting desired waves. First, due to inputting of desired waves and interfering waves having significantly high levels, it is necessary to set a source voltage or a gain of LNA without exceeding the full scale of ADC by way of the determination of RSSI1. Due to inputting of further high-level signals, it is necessary to control and reduce the gain of the preceding digital IF amplifier by way of the determination of the level of desired waves with RSSI2.\nFor the above reason, it is necessary to perform the RSSI1 operation without causing saturation of ADC. However, it is necessary to apply a variable gain to the digital IF amplifier interposed between the BPF and the demodulator in the RSSI2 operation in order to keep desired waves before the demodulator within the predetermined level without determining the existence or nonexistence of interfering waves (cf. Patent Literature Document 10 in which RSSI solely operates). This results from an expectation in which that the determination of RSSI2 will be performed with the level of a desired wave on the assumption that the level of an interfering wave has been adequately reduced by way of the band-limiting BPF. The present invention needs to compare differences between RSSI1 and RSSI2 so as to estimate the existence or nonexistence of interfering waves, thus controlling the digital AGC gain.\nThe second embodiment of Patent Literature Document 8 refers to an AGC control method using both of RSSI1 and RSSI2. The second embodiment is designed to control and greatly reduce the gain of LNA to prevent AGC saturation with respect to the very low level of a desired wave and the very high level of an interfering wave, whereas it may raise the fear of degrading the reception sensitivity due to an unnecessary reduction in the level of a desired wave. To prevent this drawback, Patent Literature Document 8 discloses a correction method of preventing a rapid reduction in a control voltage applied to LNA in order to alleviate a significant reduction in the gain of LNA which is controlled via the processing result of RSSI1.\nPatent Literature Document 11 discloses the technology similar to the technology of Patent Literature Document 5. It is possible to detect the existence of interfering waves and the levels of interfering waves by comparing differences between the level of RSSI1 (interfering waves+desired waves), preceding the analog baseband part including. MIXER through ADC or the band-limiting LPF arranged in the IF system and the level of a desired wave detected with RSSI2 following the LPF. It is necessary to control the voltage of a DC-DC converter used for LNA in order to prevent saturation of LNA, occurrence of distortion, and suppression of sensitivity (i.e. a phenomenon in which the gain or the desired-wave gain is decreased below the appropriate value due to gain compression caused by interfering waves) while improving the linearity of LNA when the strength of an interfering wave is being varied over time during reception of an interfering wave and a desired wave which are intermixed together.\nUpon determining a high level of an interfering wave, it is possible to improve the saturation power of LNA (i.e. to expand the backoff) by boosting the control voltage of LNA. In contrast, upon determining a low level of an interfering wave, it is possible to prevent an increase of power consumption, which is necessary to constantly cope with interfering waves, by restoring the original control voltage of LNA.\nPatent Literature Document 11 aims at improvement of Patent Literature Document 12. Patent Literature Document 12 aims to prevent an increase of power consumption which is necessary to constantly cope with interfering waves, wherein it determines the existence or nonexistence of interfering waves by detecting a difference of RSSI at two points, and it copes with high-level interfering waves by changing the paths regarding the LNA and the analog baseband part in the direct conversion system.\nIt is possible to determine the existence or nonexistence of interfering waves based on a difference between RSSI1 representing the RF output of the LNA reflecting interfering waves and desired waves and RSSI2 regarding the output of the digital baseband part solely reflecting desired waves, transmitted through the channel filter, after ADC. In this connection, it is possible to secure a net gain by arranging the digital AGC amplifier before ADC irrespective of a gain correction when the LNA is bypassed upon detecting interfering waves above the allowable value."}
{"text": "Conventional lift cranes include a rotatable body or upper works mounting the lift boom and machinery that rotate about a vertical axis on a lower works or body. If the crane is mobile, the lower works or body is typically crawler mounted. The lifting capacity of a mobile crane is largely determined by the geometry of the base since all of the compression and tilting loads must act through and around the mobile base to the ground. The constant demands for increasing crane capacity have been partly met by larger-sized cranes having bigger lower bodies, both for more strength and to further space the fulcrum or tipping point of the crane from the counterweight effective line of action. These larger cranes have also been provided with increasing amounts of counterweight carried on the rotatable upper works which resist the overturning moment of the larger loads.\nA significant increase in crane capacity was achieved by providing a self-propelled crane with the support ring and extended boom carrier disclosed and claimed in U.S. Pat. Nos. 3,485,383; 3,878,944; and 4,196,816. In these designs, the weight of the crane and its load is transferred to the ground through a large diameter, track-like ring. As shown in these patents, and as practiced commercially for some years, the support ring is either blocked into place by timbers fitted or wedged beneath and completely around the ring or is supported by a plurality of jacks spaced around the periphery of the ring. While such cranes have increased counterweight and lift capacities, they are no longer mobile under heavy loads.\nFurther refinements in ring supported cranes are disclosed in U.S. Pat. Nos. 4,042,115; 4,103,783; 4,387,813 and 4,387,814. These patents disclose inter alia that a separate transporter mechanism may be run in and out of an otherwise stationary ring supported crane in order to move that crane between different locations or transporter mechanisms and/or idle crawlers or dollies may be installed beneath the ring under the boom foot and counterweight. Additional ring segments of even greater radius, which may also be mounted on mobile transporter mechanisms, are disclosed in U.S. Pat. Nos. 4,316,548; 4,358,021; 4,449,635 and co-pending application Ser. No. 259,932, filed May 4, 1981, now U.S. Pat. No. 4,601,402. Other arrangements for forwardly extending the boom foot or load fulcrum point are disclosed in \"TransiLift\" type cranes such as in U.S. Pat. Nos. 3,836,010; 4,170,309; 4,243,148; 4,537,317 and 4,555,032.\nConversely, attempts have been made to increase crane capacity by adding free-swinging counterweights such as disclosed in U.S. Pat. Nos. 3,202,299 and 3,209,920 or through the use of \"Sky-Horse\" type counterweight trailers such as disclosed in U.S. Pat. Nos. 3,842,984, 3,921,815 and 4,258,852 or suspended counterweight control systems as disclosed in U.S. Pat. No. 4,557,390. Reference may also be made to much earlier movable counterweight control systems for pedestal mounted cranes such as U.S. Pat. No. 524,619 wherein the counterweight position is directly dependent on the load line tension and U.S. Pat. No. 970,773 wherein the counterweight is swung out rearwardly in opposition to the forward reach of a jib.\nNone of these prior art arrangements, however, provides a fully satisfactory arrangement for increasing the lift capacity of a mobile crane while maintaining full mobility and maneuverability of the crane on its own crawler base and also permitting full swinging movement of the crane upper under both load and no load conditions without undue swinging and counterbalancing forces being created by the suspended counterweight mechanism."}
{"text": "Coronary artery disease may produce coronary lesions in the blood vessels providing blood to the heart, such as a stenosis (abnormal narrowing of a blood vessel). As a result, blood flow to the heart may be restricted. A patient suffering from coronary artery disease may experience chest pain, referred to as chronic stable angina during physical exertion or unstable angina when the patient is at rest. A more severe manifestation of disease may lead to myocardial infarction, or heart attack.\nA need exists to provide more accurate data relating to coronary lesions, e.g., size, shape, location, functional significance (e.g., whether the lesion impacts blood flow), etc. Patients suffering from chest pain and/or exhibiting symptoms of coronary artery disease may be subjected to one or more tests that may provide some indirect evidence relating to coronary lesions. For example, noninvasive tests may include electrocardiograms, biomarker evaluation from blood tests, treadmill tests, echocardiography, single positron emission computed tomography (SPECT), and positron emission tomography (PET). These noninvasive tests, however, typically do not provide a direct assessment of coronary lesions or assess blood flow rates. The noninvasive tests may provide indirect evidence of coronary lesions by looking for changes in electrical activity of the heart (e.g., using electrocardiography (ECG)), motion of the myocardium (e.g., using stress echocardiography), perfusion of the myocardium (e.g., using PET or SPECT), or metabolic changes (e.g., using biomarkers).\nFor example, anatomic data may be obtained noninvasively using coronary computed tomographic angiography (CCTA). CCTA may be used for imaging of patients with chest pain and involves using computed tomography (CT) technology to image the heart and the coronary arteries following an intravenous infusion of a contrast agent. However, CCTA also cannot provide direct information on the functional significance of coronary lesions, e.g., whether the lesions affect blood flow. In addition, since CCTA is purely a diagnostic test, it cannot be used to predict changes in coronary blood flow, pressure, or myocardial perfusion under other physiologic states, e.g., exercise, nor can it be used to predict outcomes of interventions.\nThus, patients may also require an invasive test, such as diagnostic cardiac catheterization, to visualize coronary lesions. Diagnostic cardiac catheterization may include performing conventional coronary angiography (CCA) to gather anatomic data on coronary lesions by providing a doctor with an image of the size and shape of the arteries. CCA, however, does not provide data for assessing the functional significance of coronary lesions. For example, a doctor may not be able to diagnose whether a coronary lesion is harmful without determining whether the lesion is functionally significant. Thus, CCA has led to what has been referred to as an “oculostenotic reflex” of some interventional cardiologists to insert a stent for every lesion found with CCA regardless of whether the lesion is functionally significant. As a result, CCA may lead to unnecessary operations on the patient, which may pose added risks to patients and may result in unnecessary heath care costs for patients.\nDuring diagnostic cardiac catheterization, the functional significance of a coronary lesion may be assessed invasively by measuring the fractional flow reserve (FFR) of an observed lesion. FFR is defined as the ratio of the mean blood pressure downstream of a lesion divided by the mean blood pressure upstream from the lesion, e.g., the aortic pressure, under conditions of increased coronary blood flow, e.g., induced by intravenous administration of adenosine. The blood pressures may be measured by inserting a pressure wire into the patient. Thus, the decision to treat a lesion based on the determined FFR may be made after the initial cost and risk of diagnostic cardiac catheterization has already been incurred.\nThus, a need exists for a method for assessing coronary anatomy, myocardial perfusion, and coronary artery flow noninvasively. Such a method and system may benefit cardiologists who diagnose and plan treatments for patients with suspected coronary artery disease. In addition, a need exists for a method to predict coronary artery flow and myocardial perfusion under conditions that cannot be directly measured, e.g., exercise, and to predict outcomes of medical, interventional, and surgical treatments on coronary artery blood flow and myocardial perfusion.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory only and are not restrictive of the disclosure."}
{"text": "This invention relates to combustion apparatus and, more specifically, to means for effectively cooling combustion chambers. For convenience of illustration and discussion, the invention will be described in connection with a jet engine of the gas turbine type. However, it will be appreciated that the structure is suitable for any high temperature application which requires effective film cooling of combustion apparatus.\nAircraft engines presently in operational use and those under development for future applications are designed to operate at extremely high temperatures. In order to prolong the life of combustors associated with such engines, new alloys have been developed which are highly compatible with the high temperature environment. However, it has also been found that by cooling the combustor under operating conditions the thermal fatigue life characteristic of the combustor are enhanced.\nGenerally, it is accepted practice in the art to cool combustion chambers by providing a moving film of cooling air between the inner surface of the liner and the hot gases of combustion. The film of cooling air forms a protective barrier between the liner and the hot gases and also provides for convective cooling of the liner.\nGenerally, in prior art devices the protective film is introduced into the combustion chamber from a plenum of cooling air surrounding the exterior of the combustor. This has been accomplished by providing for the introduction of cooling air through a series of apertures in an upstream portion of the liner into an annular lipped pocket. The streams of cooling fluid entering through the apertures are permitted to mix and coalesce within the pocket to form a uniform annular boundary layer of cooling air which is directed by the lip along the inner surface of the combustor liner.\nIt is well known that the lip associated with the aforementioned cooling arrangement is subject to thermal stresses which cause warpage and buckling of the lip under operating conditions. One of the approaches utilized in the past to overcome warpage and buckling of the lip has been to include in the downstream portion of the lip a series of circumferentially spaced dimples which provide localized stiffening to resist the buckling tendency induced by the thermal stresses. While the inclusion of dimples in this manner served well to overcome lip distortion, the dimples were found to create wakes in the film of cooling air discharged along the inner surface of the liner. The wakes were found to destroy the uniformity of the cooling air barrier and permit hot gas of combustion to directly contact the inner liner of the combustor thereby reducing its operating life.\nSome attempts have been made to eliminate the wakes caused by dimples disposed in combustor lips. United States Patent 3,826,082 discloses an arrangement wherein the lateral walls of the dimple converge in the downstream direction. Unlike previous dimples which diverged in the downstream direction, the arrangement taught in the referenced patent sought to direct cooling air into the area immediately downstream of the dimple hence filling the area with cooling air rather than hot gases of combustion. Incorporation of dimples with converging lateral walls into combustors has proved to be at least partially successful in reducing the deleterious effects on the liner which result from dimple wakes. However, while converging lateral walls, as taught by the above-referenced patent, serve to form an exit slot characterized by increasing width in the aft direction, such walls do not insure that the cross-sectional flow area of the exit slot is also increasing. Rather, since the height of the dimples increases in the aft direction, the cross-sectional area flow area of the exit slot decreases in the same direction. Consequently, cooling air flowing through the exit slot of decreasing flow area must accelerate and converge as it flows through successive downstream cross sections of the exit slot. As the cooling air exits the slot in the aforedescribed converging manner, wakes are formed in the boundary layer of cooling film. Said another way, wakes are formed in the boundary layer of cooling film as a direct result of the flow of cooling air out of the exit slot in an accelerated convergent manner caused by the decreasing cross-sectional flow area of the exit slot. Such convergent flow will be present even though the pair of lateral walls associated with each dimple converge toward each other and form an exit slot with increasing width in the aft direction. The present invention is directed toward providing dimple construction in a lip associated with a cooling slot disposed in a combustor liner wherein each dimple is characterized by a constant height portion so as to provide, in cooperation with convergent lateral walls, an exit slot of increasing cross-sectional flow area in the aft direction. An exit slot having an increasing cross-sectional flow area will cause cooling air to exit the slot in a divergent manner so as to flow into and fill wake areas immediately aft of each dimple."}
{"text": "Field of the Invention\nThe invention relates to an integrated circuit for processing security-relevant data having data output circuits and access control circuits. The invention also relates to a circuit configuration for supplying power to security-relevant parts of an integrated circuit.\nIntegrated circuits used in smart cards containing security-relevant data can be the target of a wide variety of attacks on the security-relevant data contained in the integrated circuits.\nPhysical attacks on smart cards can have various goals, such as reading-out (probing) of secret signals or forcing of control signals.\nTherefore, in security technologies, secret signals and control signals are conducted in mask planes that are difficult to access, and are additionally protected by a shielding layer (a so-called security layer).\nOver and above the methods of probing and forcing, however, it is also possible to isolate circuit blocks from the supply in order deliberately to generate xe2x80x9cstuck atxe2x80x9d errors on control signals and thus to cancel e.g. blockade functions.\nIn order to defend against such attacks in which access control circuits on an IC (Integrated Circuit) are deliberately rendered voltageless, the supply of such access control circuits has hitherto been routed or conducted twice (both in the aluminum plane and in the diffusion plane). In accordance with the prior art, such an attack described above could thus be warded off by routing or conducting the supply to the access control circuits in inseparable layers (planes of the IC) for example in the diffusion plane.\nRouting or conducting the supply twice has the disadvantage that a considerable amount of space on the IC is lost, since signals could otherwise be conducted in the diffusion. Conducting the power supply exclusively in the diffusion has the disadvantage that the electrical resistance of the diffusion layer is usually higher. Therefore, either voltage drops occur, or it is necessary to provide tracks of appropriate width in the diffusion, which again leads to a considerable loss of space.\nIt is accordingly an object of the invention to provide an integrated circuit for processing security-relevant data and a circuit configuration for supplying power to security-relevant parts of an integrated circuit which overcome the above-mentioned disadvantages of the heretofore-known circuits of this general type and in which, while maintaining or improving the security, the space required for a power supply of security-relevant parts of the IC is reduced wherein the power supply is additionally or exclusively conducted in the diffusion.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, in combination with an integrated circuit having security-relevant parts and access control circuits for protecting the security-relevant parts, a circuit configuration for supplying power to the security-relevant parts, the circuit configuration including:\na power supply circuitry for supplying power to the security-relevant parts; and\nthe power supply circuitry being laid out such that a power supply to the security-relevant parts is interrupted if a power supply to the access control circuits is disturbed.\nIn other words, the object of the present invention is achieved through the use of a circuit configuration for the power supply of security-relevant parts of an integrated circuit, which are protected by corresponding access control circuits, wherein the power supply of the security-relevant parts is conducted or routed in such a way that the power supply is interrupted if the power supply of the access control circuits is disturbed.\nThe object of the invention is thus achieved by virtue of the fact that a disturbance of the power supply of the access control circuits leads to a blocking of the data output circuits.\nIn this case, a particularly simple solution is possible if the power supply of the security-relevant parts is connected to the power supply of the access control circuits.\nGreater security is afforded by a solution in which the power supply of the security-relevant parts is conducted via one or more switches which open if the power supply of the access control circuits is disturbed. In this way, it is possible to prevent the forcible re-establishment of a power supply of the security-relevant parts while the power supply of the access control circuits is interrupted.\nIn this case, it is particularly preferred for an NMOS (Negative-Channel Metal Oxide Semiconductor) switch to be provided between the general power supply VDD and the power supply of the security-relevant parts, the gate of which is connected to the VDD power supply of the access control circuits via a line routed in the diffusion or in a security layer.\nEven greater security of the integrated circuits can preferably be achieved, in conjunction, naturally, with a somewhat higher outlay, through a combination of the security measures described above.\nWith the objects of the invention in view there is also provided, an integrated circuit for processing security-relevant data, including:\ndata output circuits; and\naccess control circuits operatively connected to the data output circuits such that a disturbance in a power supply to the access control circuits results in a blocking of the data output circuits.\nOne possible preferred embodiment of this solution is based on the fact that blocking signals are generated by the access control circuits, which are respectively inverse in pairs, and the data output circuits operate only when in each case both inverse blocking signals indicate cancellation of the blocking. If one of the power supplies of the access control circuit is interrupted, one of the blocking signals inevitably assumes a xe2x80x9cfalsexe2x80x9d value, as a result of which the data output is blocked.\nIn this case, it is particularly preferred for the respectively mutually associated inverse blocking signals to be conducted parallel to one another in the integrated circuit, preferably one above the other. This makes it more difficult to attack an individual blocking signal.\nFurthermore, it is preferred for the blocking signals to be conducted in the diffusion or in a security layer. Otherwise, deblocking of the data output circuits could be achieved through an attack on the blocking signals, although with some outlay.\nAnother preferred development of the invention is based on the power supply of the data output circuits being conducted in such a way that the power supply is interrupted if the power supply of the access control circuit is disturbed.\nFor this purpose, the power supply of the data output circuits may preferably be connected to the power supply of the access control circuits. This is a very simple possibility for protecting the integrated circuit against the abovementioned manipulations.\nEven greater security is afforded by the preferred solution, in which the power supply of the data output circuits is conducted via one or more switches which open if the power supply of the access control circuits is disturbed. In this way, it is also possible to avoid the situation where the power supply of the data output circuit is re-established by placing an electrically conductive needle onto corresponding regions of the IC, even though the power supply of the access control circuits is disturbed.\nIn this case, particularly preferred is a solution wherein an NMOS switch is provided between the general supply voltage VDD and the power supply of the data output circuits, the gate of which is connected to the VDD power supply of the access control circuits via a line routed in the diffusion or in a security layer.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in an integrated circuit and a circuit configuration for the power supply of an integrated circuit, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawings."}
{"text": "The present invention relates to a method and apparatus for producing plastic products. The invention is particularly useful for recycling waste plastic materials and is therefore described below with respect to this application.\nThe increase in usage of plastic materials in everyday life has created a serious disposal problem since most types of plastic materials are not readily degradable. Efforts are therefore continually being made to develop methods for recycling waste plastic materials not only to alleviate the waste-disposal and pollution problems, but also to lower the need for virgin plastic and other raw materials (e.g., wood) in producing various products and thereby to lower the cost of such products."}
{"text": "1. Technical Field\nThe present invention relates in general to a system and method for virtualization of processor resources. More particularly, the present invention relates to a system and method for virtualizing processor memory such that one or more threads may access the processor memory, regardless of a processor's state.\n2. Description of the Related Art\nComputer systems are becoming more and more complex. The computer industry typically doubles the performance of a computer system every 18 months (e.g. personal computer, PDA, gaming console). In order for the computer industry to accomplish this task, the semiconductor industry produces integrated circuits that double in performance every 18 months. A computer system uses integrated circuits for particular functions based upon the integrated circuits' architecture. Two fundamental architectures are 1) microprocessor-based and 2) digital signal processor-based.\nAn integrated circuit with a microprocessor-based architecture is typically used to handle control operations whereas an integrated circuit with a digital signal processor-based architecture is typically designed to handle signal-processing manipulations (i.e. mathematical operations). As technology evolves, the computer industry and the semiconductor industry realize the importance of using both architectures, or processor types, in a computer system design.\nMany computer systems use a multi-processor architecture in order to provide a substantial amount of processing power while attempting to support a wide range of software applications. A challenge found, however, is that each of these processors includes dedicated internal memory that is not accessible by other processors, such as a local storage area and register files. Another challenge found is that if a particular processor is unavailable, other processors may not access data from the unavailable processor.\nFurthermore, a processor may run multiple threads that access the same memory space. A challenge found is that when a particular thread is accessing a particular memory space, the other threads must wait until the particular thread is complete before the other threads are able to access the memory space.\nWhat is needed, therefore, is a system and method to virtualize processor memory resources such that the resources are concurrently accessible by one or more threads, regardless of a processor's state."}
{"text": "Conventionally, the World Wide Web has used HTML (Hyper-Text Markup Language) to encode documents (e.g., web pages). HTML was primarily intended for human consumption and hence has limitations with respect to the ability to perform more complex functions involving web documents, which may be static or dynamic. While the content of static web documents are fixed, the content of dynamic web documents may change. In this fashion, the web document may meet a user request for information, may provide a changing graphical display, etc. Conventionally, creating a dynamic web document involves combining stylesheets and scripts to create such web documents, which change in response to user interactions.\nConventional methods to generate dynamic web pages include using a JSP (Java Server Page) or using an ASP (Active Server Page). In the JSP case, a file containing a combination of JSP tags, HTML code and Java code is transformed into Java source code, compiled, and the resulting executable is run to generate a customized web page. The JSP file is transformed and compiled the first time it is accessed and also whenever the JSP source page is newer than the compiled version. The ASP technology is analogous to the JSP technology but uses Visual Basic as a primary server side scripting language.\nUnfortunately, the above-mentioned conventional methods have several limitations. One such limitation relates to the ability to save the state associated with a web-page the user was accessing, such that the user may leave and come back to the web-page to pick up where the user left off. For example, a dynamic web-page may be designed to walk a user through of a number of steps to complete a process. A user accessing such a web-page may go off-line and come back to finish later. While some conventional systems may allow for this, the technique is often complex because the state associated with where the user is in the process must be linked to the copy of the code that was used to create the dynamic web page. Thus, multiple items must be saved (e.g., the identity of the page and the version of the code used to access it) and they must be linked together.\nSome conventional methods also face difficulties when updating the code that accesses web pages. For example, a single piece of code (e.g., a JSP) is used to process incoming requests for web pages for multiple users. From time to time, this code must be re-compiled. Using the example above in which a user leaves a process and wishes to come back later, the code may have been re-compiled between the time the user left and the time the user returned. The newly compiled code may not allow the user to complete the process properly and saving the old code is logistically unattractive. Hence, problems of mixing old and new code arise.\nAn additional problem with such conventional methods is that they are not well-suited to providing machine-readable descriptions. For example, while the markup tags in HTML define how the content is to be formatted, they do not describe the content itself. Recently, XML (Extensible Markup Language) has been used to encode web-pages. In particular, with the introduction of higher-level representations such as RDF (Resource Description Format), it is practical to provide web pages that contain machine-readable descriptions. In this evolution of the World Wide Web to the Semantic Web, it is suggested that web pages may become a store of data to be mined by autonomous software agents.\nSoftware agents are software routines that wait in the background and perform an action when a specified event occurs. Software agents may be autonomous, acting on behalf of the user. Some agents are intelligent, learning and adapting to environmental conditions. Software agents may perform a number of other tasks such as, information retrieval. In this example, a user may send a software agent to gather information on a specified topic while the user is off-line. When the user returns, the information is waiting. As another example, a software agent may function as a broker, seeking out prices for a specified product and returning a list to the users. As yet another example, software agents may transmit a summary file on the first day of the month or monitor incoming data and alert the user when a certain transaction has arrived.\nOne conventional method of coding agents is to describe the agent using a custom programming language with an XML syntax. While this may be effective for manipulating XML expressions, it requires that a new language be learned.\nOther conventional systems may not require that a new language be learned to code the agent; however, the agent itself is not directly accessible. Consequently, debugging such an agent is complicated.\nTherefore, one problem with conventional web access methods is the difficulty of saving the state associated with a web page. Another problem is the complexities involved in accessing a saved web page when the code that is used to access the web page is updated. An additional problem is the difficulty of providing web pages with machine-readable descriptions. Still other problems involve coding and debugging software routines such as agents."}
{"text": "1. Field of the Invention\nThe present invention is directed to a thrust plate assembly, especially for a motor vehicle friction clutch with automatic wear compensation. Such an assembly comprises a housing which can be fixed or is fixed to a flywheel for joint rotation about an axis of rotation. A pressure plate is arranged in the housing so as to be fixed with respect to rotation relative to it and axially displaceable relative to it. An energy accumulator, preferably a diaphragm spring, is supported at the housing on one side and is supported at the pressure plate on the other side and presses the pressure plate in the direction of a side of the housing provided for connection with the flywheel. A wear adjustment device is arranged in the support path of the energy accumulator between the energy accumulator and a component of the housing and pressure plate. The wear adjustment device has at least one adjustment element which is displaceable for purposes of wear adjustment and is pretensioned in a wear adjusting direction. At least one play sensor arrangement for detecting wear in friction facings or friction linings of a clutch disk can be clamped or is clamped between the pressure plate and the flywheel. The play sensor arrangement comprises: a locking/detection element which is arranged on one component and which is pretensioned against the wear adjustment device with a locking portion and acts upon this wear adjustment device in order to prevent a movement of the at least one wear adjustment element in the wear adjusting direction and which interacts or can be made to interact by a detection portion with another component or assembly for detection of wear. The latter component or assembly is displaceable with respect to the first component when wear occurs. The locking/detection element can be brought into a position for releasing the at least one adjustment element for movement in the wear adjusting direction when wear occurs by means of the interaction with the other component or assembly.\n2. Discussion of the Prior Art\nA thrust plate arrangement constructed in the manner described above is known from the prior art, wherein the play sensor arrangement is constructed as follows: An axial through-opening is provided in the pressure plate and is penetrated by a pin-like detection portion of a play sensor. The pin-like detection portion is easily tiltable in the through-opening. A leaf spring element is fixedly arranged at one end of this pin-like detection portion, this end being situated at a distance from the flywheel, wherein the leaf spring element extends in the direction of the wear adjustment device so that the wear adjustment device is clamped between the pressure plate and the leaf spring element. Due to the support of the leaf spring element at the wear adjustment device, the pin-like detection portion is tilted in its through-opening due to the spring elasticity of the leaf spring element, and is accordingly pretensioned in a friction clamping fit. If wear occurs in friction linings of a clutch disk which are located between the pressure plate and flywheel, the pressure plate moves in the direction of the flywheel until the pin-like detection portion stands up at the flywheel. In so doing, this detection portion is displaced in its through-opening axially with respect to the pressure plate against the pretensioning action and friction clamping fit. This means that its end by which it is connected with the leaf spring element is pushed away from the pressure plate so that the pretensioning force by which the leaf spring element acts upon the wear adjustment device decreases. When the force of the diaphragm spring exerted on the wear adjustment device is reduced or released in a subsequent disengagement process, the at least one adjustment element can move in the wear adjusting direction accompanied by increasing tensioning of the leaf spring element until the pretensioning force by which the at least one wear adjustment element is pretensioned in the wear adjustment position and the force applied by the leaf spring element balance one another, and a further movement of the at least one wear adjustment element is blocked in the wear adjusting direction by the leaf spring. This means that in this arrangement, the play sensor formed of the pin-like detection portion and the leaf spring element is displaced to an increasing extent with respect to the pressure plate and there takes place a successive relaxing of the leaf spring element (when wear is detected) and tensioning of the leaf spring element (when wear is compensated).\nIn order to carry out wear compensation in an arrangement of this kind to an extent which exactly corresponds to the wear of the friction linings detected by the play sensor, the spring elasticity or spring constant of the leaf spring element must be provided in a highly precise manner. If the leaf spring element is somewhat too soft, there is a risk that overcompensation of wear can take place; if the leaf spring element is too stiff, there is a risk that the wear will not be sufficiently compensated and that the interaction of the different components cannot take place in the provided manner."}
{"text": "In distributed storage systems, each data block may be stored as multiple replicates and when the data block is updated, the same update needs to be implemented among the replicates as well. The update among replicates can be implemented by a coordinator. In a master-slave replication system, the coordinator can be a master server and update process can be: the coordinator sequences updates data blocks, assigns new data versions, generates update requests, and sends the requests to slave servers. The update request can comprise two parts: data update and control information update. The control information update includes the new data version and metadata of the data block, such as a data block identification, the original version of the data block, the new version of the data block, and metadata information for other slave servers. Each slave server receives and processes the update request and then replies to the coordinator. The coordinator determines, according to a number of normal responses from the slave servers, the number of update requests sent, and a data consistency protocol, whether the data update is successful. If yes, it will modify the data block to the new version.\nIn the process of updating the data, the new version is generated based on the current version of the data block. Thus, the version monotonically increases and only one latest version is included in the update requests. If, instead, the slave server receives update requests including multiple new versions, the complexity will increase for maintaining consistency in updating the data. Therefore, in the update process, if the coordinator receives multiple update requests to the same data block, the request has to be queued and only after completing a current update request, it will process the next.\nIn order to achieve consistency in updating replicates of the data block, update requests of the same data block are sequenced, i.e., in serial execution of the data update requests. The data update process includes: network transmission, logic processing of the request, and write disk IO parts. At present, the disk can be solid state discs (SSD) and cluster write performance is relatively high. But carrying, in an update request, multiple data updates as batch processing can have the data update requests processed at once, reducing delay in the logic processing. The main bottleneck of data updates is in transmission network: the batch processing contains multiple updates and the amount of data transmission is large, causing a significant transmission delay."}
{"text": "Society often requires the filling out of forms, such as a government compliance form, a tax return, a loan application, or a job application. The government (both state and federal) has thousands of forms that individuals and/or companies are expected to fill out for a specific request or on a regular basis.\nA form is a list of questions organized into relevant groups of data fields that a user needs to answer. A form may have a predefined rigid structure to help the repeated filling out and the further processing and may include a field for a signature. Forms may be in paper or electronic format.\nPaper forms are often stored, transferred, and processed mostly by manual labor, that is slow and of low efficiency. The paper forms are inflexible, can be understood only via reading lengthy instructions. The instructions often refer to data that is only partially available. The instructions often use language (i.e. in government dialect) that is only partially understandable for the layperson and may not be explained. This may lead to misunderstanding and to an elaborate error-fixing process after the form is completed that necessitates time consuming, additional processing. A common practice is to hire out an expensive specialist who will fill out the forms on behalf of the individual or the company.\nAn example of a common and complex paper form is the individual tax return (Form 1040). This form contains about 150 fields (without the attachments); but an average user will need to fill out only about 60 of those fields, of which 10 are aggregated fields that could be derived from the already given answers (e.g. subtract line 33 from line 22). The instructions for the tax return comprise 32 pages.\nNot all the questions on a form are relevant to the user's situation. The paper forms often have an embedded logic that the user needs to be able to follow in order to make decisions about how to fill out the form. Typically, the more complex the form is, the less fields have to be filled out based on the user's situation. Yet the user needs to process the whole form in order to understand which fields need to be filled out and which fields can be left empty.\nThe paper forms are independent and do not know about one another. Even if the same identification data (e.g. name and address of the company, phone numbers, tax id, etc.) has to be put into each form, it has to be repeatedly put into each form separately.\nThe currently available electronic form filling solutions basically implement the paper form structure (e.g. Microsoft Word or Adobe Acrobat and other Adobe products). The decision about what fields to fill in are still based on the user's following the complicated instructions. The user still needs to go through each field and determine whether the information is needed or not. The identification data for the same user (e.g. name, address, etc.) still needs to be added to each form separately. Though these documents are often called intelligent or smart documents, the intelligence is limited to localized help, error checking, or calculations in certain data fields. The data collected in one form is not accessible to another form.\nA Microsoft Word or an Adobe PDF document can be protected by a digital signature. In the usual process of form filling, the user fills out the form, signs the document, and forwards it to the agency requiring the form. Once a document is digitally signed it cannot be converted to a programmatically accessible format from which relevant data can be extracted. The receiving agency cannot use the digitally signed document to extract e.g. statistical data via a program without further manual processing (e.g. printing, scanning, converting, and error checking again). Thus the existing solution (i.e., combining the filling out a form in Microsoft Word or Adobe Acrobat and signing it digitally) will not produce easily extractable raw data.\nSome of the currently available electronic form filling solutions (e.g. tax return generator programs like TurboTax) are able to control the questioning and thus interpret the instructions to help the user. These applications are hard coded as computer programs. If the form is changed, the corresponding program needs to be changed as well. Creating an electronic form involves the writing of a new computer program. These processes are also time consuming and error prone. Yet, a form represents knowledge and it is the nature of the knowledge to change from time to time.\nThe internet provides a convenient framework for filling out online forms. The current electronic form filling methods (e.g. those provided for the public by the government offices like the Department of Labor, the Patent Office, the Securities and Exchange, or the Internal Revenue Service) are still relegated to manually signed paper documents as the final result. Such paper documents are scanned in as images, but the images are not searchable. Though the paper images can be transformed to text file through OCR scanning, this technology requires a human final check that is still error prone and time consuming. For this reason, because of the missing electronic equivalent of a signed paper document, the paperless office is still not working on the society level.\nConsequently, and for the above delineated reasons, there exists a need for a faster, simpler way of filling out forms (both paper and online forms), creating legally binding documents that are completely electronic and not paper-based, and still permit the further processing of the accumulated data."}
{"text": "The virtual explosion in growth of multi-national companies has led to a dramatic increase in international telecommunications traffic. A significant, and ever increasing part of such traffic is devoted to conference calls, i.e., calls involving three or more parties. The manner in which such conference calls are arranged or initiated depends, among other factors, upon the type of terminal equipment available to the participants. Thus, if one of the participants has a two-line (or multi-line) capability, that party can place one call, put it on hold, place a second (or subsequent) call, and then bridge the two (or more) calls. If two or more participants have the same capability, then one participant can make all of the calls, or several participants can add one or more additional individuals to the conference, thereby in effect sharing the responsibility for initiating the conference among the parties.\nThe choice of which approach to take is not today made with any degree of consistency nor with any consideration of the advantages that may be obtained if one conference initiation approach is selected over another. In particular, it may turn out that the cost of making a call from point A to point B is less than the cost of making the same call from point B to point A, simply because of the time zone differences at those locations and the fact that the cost of making a call is time sensitive (e.g., evening or night rates are typically cheaper than day rates). \"Direction dependent\" rates are almost always encountered when international long distance calls are involved. It may also turn out that more (or better quality) circuits are available for telecommunications traffic destined for a foreign country as opposed to the circuits connecting the same endpoints but originating from locations outside of this country.\nThe situation described above is also true with respect to conference calls made using a conference facility (such as the Alliance Conference Bridge Service available from AT&T) that is located within the telecommunications network and controlled by the telecommunications provider, rather than by using customer premise equipment. In instances in which network based conference bridges are used, there are nevertheless choices that should be made to improve economy and efficiency: which participant should set up the call, which bridge location should be selected as a dominant location, what time the call should be placed, and so on."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical system of an optical pickup in an optical information recording and reproducing apparatus which records and reproduces information for optical discs having different corresponding wavelengths. More particularly, the present invention relates to an optical information recording and reproducing apparatus which allows compatibility for a plurality of optical recording mediums using laser light sources of different wavelengths, to an optical pickup, to an objective lens module, and to a diffractive optical element.\n2. Description of Related Art\nAs an optical information recording and reproducing apparatus, an optical disc apparatus is known in which recorded information can be read from an optical recording medium, that is, an optical disc, such as digital versatile disc (hereinafter, referred to as DVD), compact disc (hereinafter, referred to as CD), or the like.\nA compatible optical disc apparatus is known in which recorded information can be read from DVD and CD. As for DVD, the substrate thickness is 0.6 mm, the corresponding wavelength is in a range of 635 nm to 655 nm, and the numerical aperture (NA) of an objective lens is about 0.6. As for CD, the substrate thickness is 1.2 mm, the corresponding wavelength is in a range of 760 to 800 nm, and the numerical aperture of an objective lens is about 0.45. In the compatible optical disc apparatus, there is a case in which a laser light source having a wavelength λDVD in the vicinity of the wavelength 660 nm for DVD and a laser light source having a wavelength λCD in the vicinity of the wavelength 780 nm for CD are mounted.\nFor example, a technology is suggested in which an optical pickup device for allowing information to be recorded and reproduced for information recording mediums having different substrate thicknesses for DVD/CD, and an objective lens and an optical element used for the optical pickup device are provided (JP-A-2001-235676). The optical pickup device is suggested in which the objective lens having diffractive orbicular zones is used for the optical pickup device, such that, with an outside light flux of a predetermined numerical aperture in a use state of a small numerical aperture as a flare, recording and reproducing of information are performed for various information recording mediums having different thicknesses. The objective lens having such diffractive orbicular zones includes a diffraction surface having the diffractive orbicular zones. Here, when a function of an optical path difference of the diffraction surface is Φ(h) (where h is a distance from an optical axis), dΦ(h)/dh is a discontinuous or substantially discontinuous function at a place of a predetermined distance h.\nOn the other hand, as for blue-ray disc (hereinafter, referred to as BD), the thickness of a transmissive protection layer (which corresponds to the thickness of a transparent substrate of DVD or the like) is 0.1 mm, the corresponding wavelength is 408 nm, and the numerical aperture of an objective lens is about 0.85. Accordingly, in a BD/DVD/CD compatible optical disc apparatus, a laser light source which emits laser light of λBD in the vicinity of the wavelength 408 nm, that is, an optical system, needs to be mounted, in addition to the configuration of the above-described compatible optical disc apparatus. Further, since the optical discs of BD, DVD, and CD have different thicknesses, a unit for correcting three kinds of different spherical aberrations needs to be provided. In addition, since all of them have different numerical apertures, a corresponding unit also needs to be provided. However, in JP-A-2001-235676 described above, the specified descriptions of these units are not given. That is, it is difficult to realize compatibility of three or more kinds of recording mediums having different light source wavelengths, numerical apertures (effective diameters), optical disc thicknesses (the thickness of a transmissive protection layer), such as BD, DVD, CD, and the like by use of a single objective lens according to the related art.\nIn order to realize an optical pickup for a compatible apparatus, a method is suggested in which an objective lens exclusively used for BD and a DVD/CD compatible objective lens are used, and are switched according to wavelengths. In this case, however, since two objective lenses are used, a complex lens switching mechanism needs to be provided, which causes a problem in that manufacturing costs are increased. In addition, since an actuator is made large, it is disadvantageous to reduce the size of the apparatus. Further, a method may be considered in which an objective lens and a collimator lens are incorporated, but, since the collimator is fixed with respect to the objective lens, it may be difficult to maintain performance at the time of movement of the objective lens.\nIn any cases, if a plurality of light sources are used, and an optical system of exclusive prism, lens, and the like is configured in order to ensure compatibility of BD, DVD, and CD, an optical pickup or an overall optical head is complicated, and tends to have the large size."}
{"text": "This invention relates to methods and apparatus for controlling the direction and/or magnitude of warpage in molded plastic parts through strategic positioning of the non-homogeneous melt conditions across a stream of a laminar flowing fluid to a desirable circumferential position. This may be used in combination with more conventional process variables. The invention is useful in flow channels generally that flow a stream of laminar flowing material, such as thermoplastic or thermosetting plastics. The invention is particularly suitable for solidifying or non-solidifying runners, such as cold-runner or hot-runner injection molding machines that flow thermoplastic or thermosetting melt into a single or multiple cavity mold. The invention is also applicable to extrusion dies in which the melt conditions of the plastic can be strategically repositioned to achieve a desirable output condition from the flow channel to impart a desired material property to the flowing melt, such as to control a magnitude and/or direction of plastic part warpage.\nThermosets require heat to transition from a fluid to a solid state (the heat induces a chemical reaction) whereas thermoplastics must be cooled from a hot molten state to solidify. This is not a chemical reaction as found with thermosetting materials, but rather a phase change from liquid to solid. Thermosets are injected into a mold (via an injection molding machine or with use of a “transfer molding” process).\nWith thermoplastics the mold is cooled so that the plastic will solidify. A cold runner mold will also cool the runner after mold filling and the melt in the runner will solidify and must be removed every molding cycle. A hot runner will allow the runner material to remain molten during the entire molding cycle.\nWith thermosetting materials, the process is somewhat opposite to thermoplastics. A heated mold is used to allow the material to solidify. During injection molding or transfer molding, a fluid material is injected into a heated mold. The mold heats the material and initiates a chemical reaction causing the material to cross link and solidify. Normally the runner travels along the parting line similar to a cold runner thermoplastic mold. However, the runner is hot and the runner material solidifies with the molded parts and must be removed during every molding cycle. A cold runner system allows the material to remain fluid much like a hot runner used in thermoplastic molding.\nWarpage of plastic parts is a result of variations in shrinkage within the part as it is being formed. Sources of such warpage of molded plastic parts have previously been poorly understood. These variations in shrinkage have generally been attributed to side to side variations in mold temperature, anisotropic shrinkage variations resulting from flow induced polymer and filler orientation, and global shrinkage variations (shrinkage variations between regions of a part) resulting from differences in wall thickness, mold temperature, melt temperature and melt pressure. Accordingly, when warpage in a particular mold design was discovered, attempts to correct the warpage typically involved modification to the melt temperature, mold temperature, fill rates, or an adjustment in pack pressure or pack time, or modifications to part geometry or gate locations."}
{"text": "1. Field of the Invention\nThis invention relates to a cobalt platinum (CoPt) deposited magnetic film, and in particular, to a CoPt deposited magnetic film of thickness from 3,000 to 10,000 angstroms having a coercivity from 1,300-2,000 oersteds, and to the method of depositing such a magnetic film.\n2. Description Relative to the Prior Art\nFollowing the practice of the crystallographic art, a crystal plane is identified by indices enclosed in parentheses, e.g. (10.1), and direction in a crystal is represented by indices enclosed in brackets, e.g. [00.2]."}
{"text": "Imaging devices capable of printing images upon paper and other media are becoming increasingly popular and used in many applications including color reproduction. For example, laser printers, ink jet printers, and digital printing presses are but a few examples of imaging devices in wide use today for black and white or color imaging.\nDigital printing presses are relatively new compared with other printing technologies and may be used in place of other printing arrangements, such as analog printing presses. In one imaging example utilizing a press, a plurality of copies of the same image may be reproduced in relatively high volumes (e.g., printing business cards, catalogs, publications, etc.).\nSome imaging devices are susceptible to shifting of colors over time due to process drift in an imaging device. Some shift may be attributed to dot gain or shifts in dot gain of an imaging device where dot gain is the relationship of a printed dot area divided by a digital dot area (corresponding to an area of a dot intended to be printed).\nAt least some embodiments of the disclosure describe methods and apparatus which provide improved image formation upon media."}
{"text": "Many computing devices configured for telecommunications (“terminals”), such as smartphones, are capable of communicating via various types of networks. For example, cellular and other portable terminals may connect with circuit-switched (CS) networks such as the Global System for Mobile Communications (GSM) or more recent packet-switched (PS) networks such as Long Term Evolution (LTE) or other Internet Protocol (IP)-based networks. Network operators often maintain legacy network devices in their networks to support older protocols or communication techniques, such as CS voice calls."}
{"text": "Data parallel processing is a form of computing parallelization across multiple processing units. Data parallel processing is often used to perform workloads that can be broken up and concurrently executed by multiple processing units. For example, to execute a parallel loop routine having 1000 iterations, four different processing units may each be configured to perform 250 different iterations (or different sub ranges) of the parallel loop routine. In the context of data parallel processing, a task can represent an abstraction of sequential computational work that may have a dynamic working size that is typically determined at runtime. For example, a processor can execute a task that processes a sub-range of a parallel loop routine. In some cases, a task may be created by a thread running on a first processor and dispatched to be processed via another thread on a second processor.\nDifferent tasks may be assigned (or offloaded to) various processing units of a multi-core or multi-processor computing device (e.g., a heterogeneous system-on-chip (SOC)). Typically, a task-based runtime system (or task scheduler) determines to which processing unit a task may be assigned. For example, a scheduler can launch a set of concurrently-executing tasks on a plurality of processing units, each unit differently able to perform operations on data for a parallel loop routine.\nMulti-processor (or multi-core) systems are often configured to implement data parallelism techniques to provide responsive and high performance software. For example, with data parallel processing capabilities, a multi-core device commonly launches a number of dynamic tasks on different processing units in order to achieve load balancing. Parallel workloads may get unbalanced between various processing units. For example, while multiple processing units may initially get equal sub-ranges of a parallel loop routine, imbalance in execution time may occur. Imbalances in workloads may occur for many reasons, such as that the amount of work per work item is not constant (e.g., some work items may require less work than other work items, etc.); the capabilities of heterogeneous processing units may differ (e.g., big.LITTLE CPU cores, CPU vs. GPU, etc.); rising temperature may throttle frequencies on some processing units more than others, particularly if the heat dissipation is not uniform (as is commonly the case); and other loads may cause some processors to lag more (e.g., loads from other applications, the servicing of system interrupts, and/or the effects of OS scheduling, etc.).\nTo improve performance while conducting data parallel processing, multi-processor systems may employ work-stealing techniques in which tasks or processing units can be configured to opportunistically take and execute work items originally assigned to other tasks or processing units. For example, when a first processing unit or task finishes an assigned subrange of a shared workload (e.g., a parallel loop routine), the first processing unit may steal additional sub-ranges of work from other processing units/tasks that are still busy processing respective assignments of the shared workload. As load imbalances may often occur, work-stealing operations may allow dynamic load-balancing that improves the utilization of system processing resources and reduces the time to complete parallel work."}
{"text": "The present invention relates to microelectronic packaging of semiconductor chips and, more specifically, to IC flip chip assemblies designed to reduce the structural damage to C4 interconnections due to thermal stress and the CTE mismatch of the chip and the packaging material.\nAdvances in microelectronics technology tend to develop chips that occupy less physical space while performing more electronic functions. Conventionally, each chip is packaged for use in housings that protect the chip from its environment and provide input/output communication between the chip and external circuitry through sockets or solder connections to a circuit board or the like. Miniaturization results in generating more heat in less physical space, with less structure for transferring heat from the package.\nThe heat of concern is derived from wiring resistance and active components switching. The temperature of the chip and substrate rises each time the device is turned on and falls each time the device is turned off.\nAs the chip and the substrate ordinarily are formed from different materials having different coefficients of thermal expansion (CTE), the chip and substrate tend to expand and contract by different amounts, a phenomenon known as CTE mismatch. This causes the electrical contacts on the chip to move relative to the electrical contact pads on the substrate as the temperature of the chip and substrate changes. This relative movement deforms the electrical interconnections between the chip and printed wiring board (PWB) and places them under mechanical stress. These stresses are applied repeatedly with repeated operation of the device, and can cause fatigue of the electrical interconnections. This is especially true for the solder ball of the controlled collapse chip connection, also known as xe2x80x9cC4xe2x80x9d, connections. It is therefore important to mitigate the substantial stress caused by thermal cycling as temperatures within the device change during operation.\nOne type of semiconductor chip package includes one or more semiconductor chips mounted on a circuitized surface of a substrate (e.g., a ceramic substrate or a plastic composite substrate). Such a semiconductor chip package is usually intended for mounting on a printed circuit card or board. In the case of a ball grid array (BGA) package, the chip carrier includes a second circuitized surface opposite the surface to which the chip is attached. This, in turn, is connected to the printed circuit card or board.\nChip carriers of this type provide a relatively high density of chip connections and are readily achieved by mounting one or more semiconductor chips on the circuitized surface of a chip carrier substrate in the so-called xe2x80x9cflip chipxe2x80x9d configuration.\nFlip chip bonding is described by Charles G. Woychik and Richard C. Senger, xe2x80x9cJoining Materials and Processes in Electronic Packagingxe2x80x9d in PRINCIPLES OF ELECTRONIC PACKAGING, by Donald P. Seraphim, Ronald Lasky and Che-Yu Li, McGraw-Hill Book Company, New York, N.Y. (1988), at pages 577 to 619; and by Nicholas G. Koopman, Timothy C. Reiley, and Paul A. Totta, xe2x80x9cChip-To-Package Interconnectionsxe2x80x9d in MICROELECTRONIC PACKAGING HANDBOOK, by Rao R. Tummala and Eugene Rymaszewski, Van Nostrand Reinhold, New York, N.Y. (1988), at pages 361 to 453.\nFlip chips are small semiconductor dies having terminations all on one side of the entire face of the die in the form of a pattern of solder pads or bump contacts. These solder bumps are deposited on solder wettable terminals on the chip. Typically, the surface of the chip has been passivated or otherwise treated. In this way the use of a flip chip package allows full population area arrays of I/O. The flip chip derives its name from the practice of flipping or turning the chip over after manufacture, prior to attaching the chip to a matching substrate.\nAs described by Seraphim et al. and Tummala et al., an electronic circuit contains many individual electronic circuit components: thousands or even millions of individual resistors, capacitors, inductors, diodes, and transistors. These individual circuit components are interconnected to form the circuits. The individual circuits are interconnected to form functional units. Power and signal distribution are performed through these interconnections. The individual functional units require mechanical support and structural protection. The electrical circuits require electrical energy to function and the removal of thermal energy to remain functional. Microelectronic packages such as chips, modules, circuit cards, circuit boards, and combinations thereof are used to protect, house, cool, and interconnect circuit components and circuits.\nIn the flip chip configuration, the chip or chips are mounted active-side-down on solderable metal pads on the substrate using solder balls, a C4 connection, a gold bump, or a conductive epoxy.\nControlled collapse chip connection in flip chip technology has been successfully used for about 30 years for interconnecting high I/O count and area array solder bumps on the silicon chips to the base ceramic chip carriers (e.g., alumina carriers). In the C4 process, as distinguished from the flip chip process, the solder wettable terminals on the chip are surrounded by ball limiting metallurgy (BLM), and the matching footprint of solder wettable terminals on the card is surrounded by a solder mask. These structures act to limit the flow of molten solder during reflow.\nBonding can be used in an unpackaged configuration known in the art as direct chip attach (DCA): the direct connection of a chip to a card or board without an intermediate layer of packaging. The combination of a chip mounted on an intermediate carrier is usually described as the xe2x80x9cfirst level packagexe2x80x9d. DCA can be a lower cost method of connecting a chip to a card, since the first level package is thereby eliminated.\nFor direct chip attach, individual IC chips are mounted on the cards or boards. The space between the mounted chip and the card or board is then filled with an epoxy resin. By this expedient, the standoff between the IC chip and the card or board is encapsulated with epoxy.\nIf a polymeric dielectric card or board is employed, the DCA process requires low temperature solder metallurgy. Moreover, direct chip attach, when used with an underfill, increases the fatigue resistance of the C4 solder interconnections to thermal cycling, acts as an alpha emission barrier to MOSFET memory chips, is a parallel thermal path for heat dissipation, and provides physical protection to the chips and C4 solder interconnections.\nHowever, one problem encountered with the combination of DCA and C4 bonding is the difficulty of reworking the encapsulated package. In order to improve rework and to accommodate the CTE mismatches between the chip and the PWB, many prior art proposals have been developed to connect integrated circuit chips to printed wiring boards via an intermediate element. Often, chip carriers are interposed between the chip and the circuit board; the CTE of the chip carrier is itself chosen as some intermediate value to provide a reasonable match to both the chip and to the printed circuit board. The very large difference in CTE between the silicon device and the printed circuit board generally requires some intermediate device carrier. One such type of interconnection mounts the integrated circuit chip on a ceramic chip carrier or module, which module is mounted on a circuit board. One or more chips may be mounted on each device carrier or module, and one or more modules may be mounted on any given circuit board. In a particularly well known type of configuration, the integrated circuit chip is mounted onto a ceramic module by flip chip bonding wherein the I/O pads on the face of the chip are bonded to corresponding pads on the module. Such connections are formed by solder bumps or solder balls normally using solder reflow techniques. It is these connections that are referred to as C4 connections.\nMost conventional single and multiple chip packages are typically constructed from thick, mechanically robust, dielectric materials, such as ceramics (e.g., alumina, aluminum nitride, beryllium oxide, cordierite, and mullite) and reinforced organic laminates (e.g., epoxies with woven glass, polyimides with woven glass, and cyanate ester with woven glass). In some cases, materials are combined to produce certain improved properties. For example, a package may have a ceramic base with one or several thin films of polyimides or benzocyclobutane (BCB) disposed thereupon.\nIn an attempt to overcome the problem of thermal mismatch between the chip carrier and the circuit board it has been proposed to fashion the chip carrier from a material similar to that of the circuit board. Such techniques are described in IBM Technical Disclosure Bulletin Vol. 33, No. 2, pages 15-16 and IBM Technical Disclosure Bulletin Vol. 10, No. 12, pages 1977-1978. However, both of these references require that the connections, at least for the signal I/O lines, be on the same side of the carrier as that to which the chip is mounted. These techniques do solve the problem of thermal mismatch between the chip carrier and the circuit board, but require peripheral I/O bonding and an additional interposer between the chip and the chip carrier. IBM Technical Disclosure Bulletin Vol. 10, No. 12 requires peripheral attachment of the chip to an interposer (carrier 2) which is bonded to the chip carrier and then attached to the card. This peripheral bonding on the chip limits the I/O that can be placed on a small chip.\nAs noted above, whether the chip is attached to a carrier or directly to the PWB, these structures are made of materials with coefficients of thermal expansion that differ from the CTE of the material of the semiconductor device: silicon. Normally the device is formed of monocrystalline silicon with a coefficient of thermal expansion of 2.5-3.0 ppm/xc2x0 C. If the substrate is formed of a ceramic material, typically alumina with a coefficient of expansion of 5.5-6.5 ppm/xc2x0 C., then the mismatch in thermal expansion is relatively small. The problem is exacerbated with less expensive substrate materials, such as fiberglass, that have a CTE of roughly 17 ppm/xc2x0 C.\nThe stress on solder bonds during operation is approximately proportional to (1) the magnitude of the temperature fluctuations, (2) the distance of an individual bond from the neutral or central point (DNP), and (3) the difference in the coefficients of expansion of the material of the semiconductor device and the substrate. Stress is also inversely proportional to the height of the solder bond (i.e., the spacing between the device and the support substrate). The seriousness of the situation is further compounded by the fact that, as the solder terminals become smaller in diameter in order to accommodate the need for greater density, the overall height decreases.\nIn order to strengthen solder joints without affecting the electrical connection, the gap is filled with a polymeric encapsulant, typically a filled polymer. The encapsulant is typically applied after the solder bumps are reflowed to bond the integrated circuit die to the printed circuit board. A polymeric precursor is dispensed onto the board adjacent the die and is drawn into the gap by capillary action. The precursor is then heated and cured to form the encapsulant. This curing can also create stresses that can be detrimental to the die.\nAn improved solder interconnection structure with increased fatigue life was disclosed in U.S. Pat. No. 4,604,644 to Beckham, et al., assigned to the present assignee, disclosure of which is hereby incorporated by reference. The Beckham patent discloses a structure for electrically joining a semiconductor device to a support substrate that has a plurality of solder connections. Each solder connection is joined to a solder wettable pad on the device and to a corresponding solder wettable pad on the support substrate. Dielectric organic material is disposed between the peripheral area of the device and the facing area of the substrate. The material surrounds at least one outer row and column of solder connections but leaves the solder connections in the central area of the device free of dielectric organic material.\nOther prior art solutions make use of an underfill material disposed between the chip and the supporting substrate in an attempt to redistribute the stress caused by CTE mismatch. Without the underfill material, this stress is wholly borne by the solder balls. The underfill material allows this stress to be more uniformly spread out over the entire surface of the chip, supporting substrate and solder balls. Examples of the use of underfill materials may be found in U.S. Pat. Nos. 5,194,930, 5,203,076, and 5,249,101.\nAfter soldering the IC to the substrate, an epoxy resin or other material is inserted into the space between the IC and the substrate and acts as a glue. In addition to being inserted into the space, surface tension produces a capillary action between the IC and the substrate that pulls the epoxy into the space. The epoxy is also pulled up along the sides of the IC by the surface tension. To make the mechanical bond of the epoxy even stronger, it is possible to roughen the surface of the substrate or the IC, chemically for instance, before applying the epoxy underfill.\nAn epoxy, after being introduced under capillary action into a space provided between the semiconductor chip and the package substrate, is cured and hardened. The hardened epoxy acts to bond the semiconductor chip to the package substrate and to protect the fragile solder connections.\nThe normal flow in a flip chip attach process is as follows: a) the die is fluxed; b) the die is placed on the substrate with bond pads on the die being aligned with bond pads on the substrate; c) solder is reflowed between bond pads on the die and substrate; d) the die is underfilled with a thermoset material; and e) the underfill material is fully cured.\nUnderfilling ensures minimum load on the interconnects and becomes the primary load bearing member between the chip and the substrate during thermal or power cycling induced due to the operation of the chip. Thermoset type materials are commonly used in the industry as underfill material. In order for the epoxy to bear much of the load it must be relatively rigid. As a result the chip and substrate are strongly coupled so that differential thermal expansion causes bending of the package when the temperature varies from that at which the epoxy was cured. In extreme cases this bending can cause cracking of the chip but, with proper design, this problem may be overcome. Without the epoxy this differential expansion causes shear in the C4 solder joints but little package bending. Warpage tends to be smaller when thicker organic substrates are used and greater as the rigidity and thickness of the substrate decreases.\nAccordingly, it is a general object of the present invention to provide novel and useful semiconductor devices wherein the foregoing problems are mitigated.\nOne specific object of the present invention is to provide a package for the semiconductor chip that minimizes stresses and strains that arise from differential thermal expansion on the chip-to-substrate or chip-to-card interconnections.\nAnother object of the present invention is to provide a package that can be readily modified to changes in chip size and configuration.\nAnother object of the present invention is to design a package that allows adhesive or underfill to be applied after the package has been assembled.\nIn accordance with the present invention, there is provided an electrical package that includes a substrate with an upper surface, having at least one electrical circuitry connecting pad. An electric device is also provided with a first modulus of elasticity. The electric device has a connection region on its lower major surface. At least one flexible connector links the connecting pad(s) to the connection region. A collar is also provided with a second modulus of elasticity. The vertical side or dimension of the collar is shorter than the electric device perimeter and is bonded to the perimeter. The horizontal side or dimension of the collar is bonded to the substrate upper surface. The adhesive material has a third modulus of elasticity. Thus, a unitary electrical package is formed that includes the electric device, the substrate and the collar. The second modulus of elasticity is at least as large as the third module of elasticity."}
{"text": "Current child-sized toy vehicles come in a variety of different models, which are typically targeted to children of different age ranges. For example, a foot-to-floor toy vehicle is a model which is typically intended for children in the age range of 12-36 months. Foot-to-floor vehicles are propelled when a child, seated on the vehicle, pushes his or her foot against the ground. A pedal toy vehicle is a model which is typically intended for children in the age range of 2-7 years. Pedal vehicles are powered by manually rotating foot pedals which are attached to a transmission unit attached to the axle of the vehicle. An electric toy vehicle is a model which is typically intended for children in the age range of 3-7 years. Electric vehicles are powered by activating an electronic transmission unit which rotates an axle of the vehicle. A gas-powered toy vehicle is a model which is typically intended for children in the age range of 6-11 years. Gas-powered toy vehicles are powered by a gas powered transmission unit which rotates an axle of the vehicle.\nEach of the above described models can be produced in a variety of different sizes, shapes, colors and body styles. For example, some toy vehicle makers have produced toy vehicles which resemble various different body styles, such as sports cars, trucks or Jeeps®. However, these toy vehicles do not allow for the interchangeability of body styles, such that one toy vehicle can be assembled, for example with sports car components and later disassembled and replaced with truck components, or assembled as a hybrid vehicle, for instance, having a sports car front end and engine, but a truck rear end and tires.\nChildren in the intended age range of these toy vehicles can have short attention spans and become bored with individual toys rather quickly. As a result, producing a toy vehicle having interchangeable body styles allows a child to use his or her own creative abilities to transform the toy vehicle according to the child's whimsical desires.\nConsequently, an improved child-sized toy vehicle is needed to address the problems of the prior art."}
{"text": "The inventions described below relate to the field inflatable liners for shoes and boots, and their method of manufacture. More specifically, the invention relates to unique removable snowboard boot liners made from Ethyl Vinyl Acetate (EVA) containing an air bladder and their method of manufacture.\nWhile removable liners present many advantages, one major disadvantage associated with removable liners is that the foot of an athlete wearing a snowboard boot with a removable liner has a tendency to slip and move within the boot. This slipping may be caused by the athlete\"\"s foot moving within the liner, by the lliner moving within the boot, or by a combination of these two phenomena. The most common result of this slipping is that the heal of the athlete lifts up from the bed of the boot. This slipping and lifting makes it more difficult for the athlete to control the snowboard, results in blisters and increases the likelihood of more serious injuries to the athlete.\nThis slipping problem has been found in the context of ski boots. Multiple solutions to this problem in the field of ski boots have been presented including several variations on the theme of using an air bladder positioned at various locations between the foot of the athlete and the exterior of the boot. Air bladders also have been used in other types of footwear for various purposes.\nThis invention is specifically concerned with snowboard boot liners made of EVA having an inflatable air bladder incorporated therein and the mass production of such liners.\nHolstine, U.S. Pat. No. 5,692,321 discloses an athletic boot such as a snowboard boot having a bladder system consisting of an upper and lower bladder in communication with each other disposed between the wearer\"\"s foot and the exterior of the boot. The upper and lower bladder system is a closed system. Thus, the overall inflation level of the bladder system is fixed and may not be readily adjusted. The system disclosed in Holstine is designed to give increased support to the ankle of a wearer when downward forces compress the lower bladder causing a corresponding inflation of the upper bladder. Holstine stresses that the disclosed boot provides the athlete with increased flexibility and range of motion of the athlete\"\"s foot when impact or operational forces are removed from the boot. Holstine does not disclose how his boot is to be produced.\nPotter, et al., U.S. Pat. No. 5,765,298 discloses an athletic shoe with an inflatable bladder present in the ankle collar The bladder has weld fines or other means incorporated therein to prevent the formation of restrictive vertical columns of pressurized gas in the medial and lateral section of the bladder. This allows increased flexibility and mobility of the wearer\"\"s ankle. Potter, et al. do not disclose how their product is to be produced.\nNishimura U.S. Pat. Nos. 3,744,159 and 3,758,964 disclose sports shoes containing an inflatable air bladder. The majority of Nishimura\"\"s disclosures focus on ski boots containing inflatable air bladders. Nishimura does not disclose the methods used to manufacture his products. Nor does Nishimura disclose the use of EVA as a liner material.\nNone of the foregoing prior art have suggested snowboard boot liners made of EVA having an inflatable air bladder incorporated therein and the mass production of such liners.\nThe present invention is directed to snowboard boot liners made of EVA having an inflatable air bladder incorporated therein and the mass production of such liners.\nThe steps involved in making a conventional snowboard boot liner from EVA are as follows:\n1. The pattern of the liner must be cut from a flat sheet of EVA;\n2. The EVA is then folded and stitched together to approximate the shape of the final liner;\n3. The stitched EVA is placed on a last;\n4. The stitched EVA on the last is heated to allow the EVA to be molded to its final shape; and\n5. The heated EVA is molded to its final shape.\nThe method outlined above is well known to those in the art. The length of time and temperature used to heat the EVA in step 4 is well known to those ordinarily skilled in the art of making snowboard boot liners from EVA. However, for purposes of illustrative example, the heating may be accomplished in a tunnel oven set at about 100xc2x0 C. The total residence time in the tunnel oven may be approximately five minutes. The exact methods used to accomplish the molding of the EVA to its final shape in step 5 are well known to those of ordinary skill in the art of making snowboard boot liners from EVA. However, for purposes of illustrative example, neoprene socks may be used.\nAir bladders in shoes generally are made by placing one sheet of suitable material on top of a second sheet of suitable material, cutting out the material in the appropriate pattern and securing the sheets together to form the boundary of the bladder and any desired internal contours such as uninflated spaces to accommodate the ankle bones. Frequently the two sheets are attached together by melting both sheets together wherever a seam is desired.\nThe method of the invention allows the air bladder to be stitched into place on the EVA liner before the EVA is stitched and molded. This allows the air bladder to be stitched when the liner is still a flat sheet of EVA. It would be extremely difficult, if not impossible, to attach the air bladder by machine stitching once the EVA has been folded, stitched, heated and molded. Prior to the invention of the method described in this patent, it was not possible to attach an air bladder made from typical materials known in the art of footwear air bladders prior to the heating and molding of the EVA. This is because the two sheets that make up the air bladder would melt together when the bladder went through the heating and molding process with the EVA-Thus, the bladder could not be inflated and would be useless. This made the mass production of EVA snowboard boot liners with inflatable bladders impractical.\nThe present method solves this problem by introducing a substance between the two sheets that make up the air bladder prior to the sealing of the seams of the bladder and partially inflating the air bladder prior to the heating and molding of the EVA liner. The introduction of this substance and partial inflation of the bladder keep the two sheets that make up the bladder separated during the heating and molding of the EVA. In addition, the disclosed design and installation of the bladder minimizes the amount of slipping of the athlete\"\"s foot within the liner and minimizes the amount of lifting of the athlete\"\"s heal from the bed of the boot."}
{"text": "The present invention relates to an positive yarn feeding device, and more particularly to an positive yarn feeding device in which the operation speed can be optionally adjusted and the yarns with different colors can be easily interchanged.\nAn positive yarn feeding device is developed for knitting colorful cloth from various colors of yarns. The early positive yarn feeding device employs electromagnetically controlled measures in the yarn feeding operation. For example, Taiwan Invention Patent Application No. 7113777 discloses a yarn feeding device for optionally feeding several yarns to a horizontal yarn knitting machine. In such device multiple sensors 7A-7D, D1/T1, D2/T2, D3/T3, D4/T4, etc. for detecting the position of the yarn guiding arms, and multiple relays CR, timers TC, TD1, TD2 and quite complicated logic control circuit are necessary to form the device.. Therefore, the cost thereof is very high and cannot be accepted in the market. Moreover in the site of the knitting machine, the fluffs always suspend in the air and attach to the device so that the logic circuit cannot work stably. Above all, once the contacts of the relays' circuit and timer are jammed with fluffs, the work thereof will fail and the whole control circuit will malfunction.\nTaiwan Utility Model Patent Application No. 7223942 discloses a positive yarn feeding device for a knitting machine which improves the above technical problem, in which when the knitting machine needs yarns in knitting operation, the tension of the yarn is increased and when the tension is greater than a preset value of the yarn feeding device, the yarn on the yarn feeding roller is pulled by a movable yarn guiding arm to contact with the frictional yarn feeding surface of the yarn feeding wheel so as to achieve the object of positive yarn feeding. Such structure is simple and practical. However, recently the knitting market tends to develop the knitting technique of multiple colors and twin yarn feeders so as to achieve various colors and figures of the product. When it is desired to produce two colors or three colors of alternately knitted cloth, two or three layers of yarn feeding wheels are necessary to be vertically piled for meeting the factory technical requirement. As shown in FIG. 9, the movable yarn guide 45 and the tension control means are disposed right under the yarn feeding wheels 17, 18 and the drive pully 16. Accordingly, when two colors or four colors of yarns are fed under the same yarn feeding wheel at the same time for knitting a cloth with colorful stripes, two or over two layers of yarn feeding wheels will be needed. This greatly increases the height of the space for the operation and causes difficulty in yarn threading and connecting processes. Moreover, when controlling or adjusting the speed, such height is always far beyond the head of an operator and thus causes great inconvenience. In addition, multiple yarns are parallelly disposed and tend to interfere with one another. In case there are three or more than three layers of yarn feeding wheels, the errors in operation will take place even more frequently. Above all, when the feeding speed of the yarn feeding wheels on the same layer cannot meet the knitting requirement and the yarn needs to be guided to the more upward or more downward layer of yarn feeding wheel, the yarn must be first torn off and then the yarn connecting and threading processes can be performed. Accordingly, the knitting changing operation will be quite time-costing and labor-costing. Besides, in such a yarn feeding operation, the main yarn feeding technique is such that whether the yarn feeding operation should be performed is decided by the effect of the frictional outer periphery of the yarn feeding wheel and the necessary tension of the yarn. However, when the yarn feeding wheel does not feed the yarn, as shown in FIG. 10, the yarn 76 still contacts with the frictional periphery of the yarn feeding wheel by at least 90 degree contacting angle. If the frictional coefficient of the yarn feeding wheel is too great, mistakes in operation will take place. The mistake which occurs most often is that when the yarn is not needed and the contacting angle between the yarn 76 and the frictional periphery of the yarn feeding wheel is considerably large, the yarn will be still fed to cause negative yarn feeding phenomenon. Therefore, the space between the stripes of the knitted cloth will be irregular or the stripes will be oblique. While the fixed yarn guiding arm 25 is adjustable as shown in FIG. 11 to minimize the contacting angle between the yarn 76 and the periphery of the yarn feeding wheel 17 or maximize the contacting angle to be over 180 degrees when positively feeding the yarn, because the yarn guiding arm 25 of the yarn feeding wheel is disposed above the yarn feeding wheel, when the number of the yarn feeding wheel is over 120, it is very difficult or even impossible to adjust the yarn feeding wheels one by one. In addition, in two color or three color knitting operation, it will be quite difficult to operate or adjust the second or third layer wheel since the height thereof will be unreachable. Therefore, the fixed arm 25 in fact is unable to meet the requirement of the present knitting market.\nMoreover, if the yarn feeding wheels have the same operation speed for multicolor knitting operation, although the wheels can be disposed on the same level, the space between the wheels will be too small to operate the yarn. This will affect the time for changing the knitting operation and the amount of the products or even indirectly cause failure of the device. If the currently most popular elastic rubber yarn or metal yarn is to be mixed with the general yarns, because the space is limited, such mixing process cannot be accomplished. Therefore, the pattern of the produced cloth will not satisfy the present market requirement. In case each yarn feeding wheel feeds two yarns at the same time barely, since the yarn guiding eye thereof has only one single eyelet, although the tension of the yarn can separate the adjacent two yarns when the yarns are fed, once the yarn feeding operation is stopped and the tension of the yarn disappears, the adjacent two or three yarns will twist with one another due to the self-twisting property thereof. Therefore, when the device is again activated, the failure thereof will inevitably occur. Therefore, the yarn feeding wheel can feed only one single yarn."}
{"text": "The present disclosure generally relates to a rearview device system, and more particularly, to a display mirror assembly having a partially reflective, partially transmissive element and a display behind the reflective element."}
{"text": "U.S. Pat. No. 4,836,767 (Schad) relates to an apparatus for producing molded plastic articles which is capable of simultaneously producing and cooling the plastic articles. The apparatus has a stationary mold half having at least one cavity, at least two mating mold portions, each having at least one core element, mounted to a movable carrier plate which aligns a first one of the mating mold portions with the stationary mold half and positions a second of the mating mold portions in a cooling position, a device for cooling the molded plastic article(s) when in the cooling position, and a device for moving the carrier plate along a first axis so that the aligned mold portion abuts the stationary mold half and the second mating mold portion simultaneously brings each plastic article(s) thereon into contact with the cooling device. The carrier plate is also rotatable about an axis parallel to the first axis to permit different ones of the mating mold portions to assume the aligned position during different molding cycles.\nU.S. Pat. No. 6,299,431 (Neter) discloses a rotary cooling station to be used in conjunction with a high output injection molding machine and a robot having a take-out plate. A high speed robot transfers warm preforms onto a separate rotary cooling station where they are retained and internally cooled by specialized cores. The preforms may also be simultaneously cooled from the outside to speed up the cooling rate and thus avoid the formation of crystallinity zones. Solutions for the retention and ejection of the cooled preforms are described. The rotary cooling station of the present invention may be used to cool molded articles made of a single material or multiple materials.\nU.S. Pat. No. 6,391,244 (Chen) discloses a take-out device for use with a machine for injection molding plastic articles such as PET preforms. The take-out device has a plurality of cooling tubes that receive hot preforms from the molding machine, carry them to a position remote from the molds of the machine for cooling, and then eject the cooled preforms onto a conveyor or other handling apparatus. The preforms are retained within the cooling tubes by vacuum pressure, but are then ejected by positive air pressure. A retaining plate spaced slightly outwardly beyond the outer ends of the cooling tubes is shiftable into a closed position in which it momentarily blocks ejection of the preforms during the application positive air pressure, yet allows them to be dislodged slightly axially outwardly from the tubes. Such slight dislodging movement is inadequate to vent the air system to atmosphere such that sufficient dislodging air pressure remains in tubes where the preforms might otherwise tend to stick and resist ejection. After the momentary delay, the plate is shifted to an open position in which all of the dislodged preforms are freed to be pushed out of the tubes by the air pressure. Preferably, the retaining plate is provided with specially shaped holes having pass-through portions that become aligned with the tubes when the plate is in its open position, and smaller diameter blocking portions that become aligned with the tubes when the plate is in its closed position. The smaller diameter blocking portions exceed the diameter of the neck of the preforms but are smaller in diameter than the flanges of the preforms such that surface areas around the blocking portions overlie the flanges to block ejection of the preforms as they undergo their dislodging movement.\nEP Pat. No. 1515829 (Unterlander) relates to a method and apparatus for cooling molded plastic articles after molding is finished. In particular, the disclosed invention relates to method and apparatus for a post mold cooling (“PMC”) device having at least two opposed faces. The method and apparatus are, according to the inventors, particularly well suited for cooling injection molded thermoplastic polyester polymer materials such as polyethylene terephthalate (“PET”) preforms."}
{"text": "Inventive concepts described herein relate to a semiconductor memory device, and more particularly, to a nonvolatile memory device and/or a control method thereof.\nSemiconductor memory devices may be characterized as volatile or nonvolatile. A volatile semiconductor memory device may perform read and write operations in a high speed, while data stored therein may be lost upon removal of power from the device. A nonvolatile semiconductor memory device may retain data stored therein even upon removal of power from the device. Accordingly, a nonvolatile semiconductor memory device may be used to store data which must be retained regardless of whether the device is powered.\nA flash memory device may be a typical nonvolatile semiconductor memory device. The flash memory device may be used as a voice and image data storing medium for information appliances such as a computer, a cellular phone, a PDA, a digital camera, a camcorder, a voice recorder, an MP3 player, a handheld PC, a game machine, a facsimile, a scanner, a printer, and the like.\nIn recent years, research has been conducted on nonvolatile memory devices where memory cells are stacked three-dimensionally to improve a degree of integration of the devices. Such a nonvolatile memory device may be referred to as a vertical NAND flash memory device or a three-dimensional nonvolatile memory device. In the three-dimensional nonvolatile memory device, word lines may be stacked in a direction perpendicular to a substrate. Cell strings may be constituted by forming pillars to penetrate the stacked word lines.\nThe pillars may be disposed in a zigzag shape to improve efficiency of space arrangement. This structure may cause irregular driving characteristics of strings. Thus, a technique of solving such a problem of the three-dimensional nonvolatile memory device is desired."}
{"text": "The present invention relates to a tape-like carrier for mounting of integrated circuit, and more particular to the tape-like carrier used for tape automated bonding in which plural inner leads are automatically bonded to a semiconductor pellet. Further the present invention relates to the tape-like carrier suitable for connection of very multiple terminals of a semiconductor pellet with electronic elements.\nHeretofore well known is a longitudinally extending tape-like carrier for mounting of integrated circuit including sets of leads secured to a flexible insulating film having longitudinally spaced apertures which are dimensioned to encompass contact regions of semiconductor pellet, said leads extending from the periphery of the every aperture toward the inside thereof (U.S. Pat. No. 3,689,991).\nFIGS. 1 to 2b show one prior illustration of such a tape-like carrier for mounting of integrated circuit, together with a semiconductor pellet having bump electrodes with leads, FIG. 1 being a plan view, FIG. 2a being a cross-sectional view taken on line IIa-IIa of FIG. 1 and FIG. 2b being a cross-sectional view taken on line IIb-IIb of FIG. 1.\nWith reference to FIGS. 1 to 2b, the conventional tape-like carrier for mounting of integrated circuit comprises a flexible insulating film 5 in the form of tape made of polyimide. The flexible insulating film 5 is provided with sprocket holes 7 at regular intervals along the longitudinal direction in both side regions thereof. Further the flexible insulating film 5 has longitudinally spaced apertures 26 which are dimensioned to encompass contact regions of a semiconductor pellet.\nLongitudinally spaced sets of leads 23 and bars 28 in the form of a frame leading to the leads 23 are formed on the flexible insulating film 5. The leads 23 extend within the periphery of a respective aperture 26 for semiconductor pellet. Between the leads 23 and the leads 23 and between the leads 23 and the bars 28, a surface of the flexible insulating film 5 is partially exposed.\nLocated in the center of the aperture 26 is a semiconductor pellet 1 having plural bump electrodes 2 formed on the contact regions of the pellet 1. The bump electrodes 2 are bonded to the corresponding leads 23 which are arranged at regular intervals so as to be coincident with the bump electrodes 2, by a heat bonding method using a bonding tool, that is, gang bonding.\nWhen bonding of very multiple inner leads 23 with very multiple bumps 2 of the semiconductor pellet is performed using the above tape-like carrier for mounting of an integrated circuit, the heat dissipation characteristic is remarkably worse so that the following inconvenience occurs: i.e., as metal regions which dissipate heat of the leads 23 conducted from the bonding tool are only the surfaces of the bars 28 and the leads 23, the efficiency of heat dissipation is low.\nTherefore, when the inner lead bonding is performed by using the above tape-like carrier for mounting the integrated circuit, thermal expansion of the leads 23 and the deformation and contraction of the flexible insulating film 5 occur, which results in the change of pitches between the leads 23. Thereby, an outer lead bonding for mounting the semiconductor pellet 1 to a substrate (not shown) or to an outer lead frame cannot be normally performed in a postprocess, which results in the faulty bonding."}
{"text": "1. Field of the Invention\nThe present invention relates to a mass spectroscopy spectrum analysis system using a mass spectrometer, and to a system for automatically determining an optimum flow of mass spectroscopy within a measurement time in order to identify the chemical structure of biopolymers, such as polypeptides or sugars, with high precision and efficiency.\n2. Background Art\nIn a general mass spectroscopy, a sample as the object of measurement is ionized, and a variety of resultant ions are delivered to a mass spectrometer for measuring the ion intensity for each mass-to-charge ratio m/z, which is the ratio of the mass number m of ion to the valence z. As a result, a mass spectrum is obtained, which consists of a peak of the measured ion intensity (ion peak) for each mass-to-charge ratio m/z value. Such a mass spectroscopic analysis of the ionized sample in a first dissociation step is called MS1. In tandem mass spectrometer, in which multiple-stage isolation is possible, an ion peak having a specific mass-to-charge ratio m/z is selected (the selected ion species is called a parent ion) from the ion peaks detected by MS1, and the thus selected ion is dissociated and broken up by collision with gas molecules or the like. The resultant dissociated ion species is then subjected to mass spectroscopy, thereby obtaining a mass spectrum in a similar manner. The n-stage dissociation of the parent ion and the subjecting of the dissociated ion species to mass spectroscopy are referred to as MSn+1. Thus, in the tandem mass spectrometer, the parent ion is dissociated in multiple stages (1, 2, . . . , n stages), and the mass number of the ion species generated in each stage is analyzed (MS2, MS3, . . . , MSn+1).\n(1) Most of the mass spectrometers capable of tandem analysis are equipped with a data-dependent function whereby, when selecting the parent ion for MS2 analysis from the ion peaks in MS1, the ion peaks are selected in decreasing intensities (such as the ion peaks in the top 10 strongest-intensities) as the parent ions, and then they are subjected to dissociation and mass spectroscopy (MS2).\n(2) The ion-trapping type mass spectrometer manufactured by Finningan is equipped with a Dynamic Exclusion function whereby, when selecting a parent ion for MS2 analysis from the ion peaks in MS1, the ion species having a mass-to-charge ratio m/z value that is designated by the user in advance is excluded from the selection as a parent ion.\n(3) Known examples relating to the determination of correspondence between a measured ion species and an ion species that has been measured include the following:\nPatent Document 1: JP Patent Publication (Kokai) No. 2001-249114 A\nPatent Document 2: JP Patent Publication (Kokai) No. 10-142196 A 1998\nIn Patent Document 1, a characteristic peak in the first-stage spectrum data and the spectrum data in the second stage of the corresponding ion species are stored in a database. In the subsequent measurements, spectrum data obtained by mass spectroscopy in the second stage of a sample as the object of measurement is compared with the second-stage spectrum data in the database in order to determine the degree of correspondence. Data components with the highest degree of correspondence is outputted as the comparison result.\nIn Patent Document 2, a measurement is continuously carried out during a multiple-stage dissociation measurement without conducting a sample injection process during measurement so that an ion intensity fluctuation due to injection between the MSn and MSn+1 data can be prevented. In this way, the need for the addition of a standard sample can be eliminated, thereby enabling an efficient quantitative analysis. The routine returns to MSn+1 or proceeds to the next MS1 measurement, depending on whether or not the data corresponds to the designated ion data that has been already collected in the MSn and MSn+1 data analysis.\nReviews of Modern Physics, Vol. 62 (1990), pp. 531-540, provides a basic description of an ion trap. A cross section of a basic configuration of the ion trap is shown in FIG. 15. The ion trap, which is a quadrupole ion trap, is made up of two end-cap electrodes and a single ring electrode. An RF voltage is applied to these electrodes such that a quadrupole electric field is formed at the center of these electrodes, thus enabling the trapping of gaseous ions three dimensionally. By continuously varying the RF voltage, the mass of the ions that are discharged can be controlled. A quadrupole pole is made up of four parallel poles. By applying a RF voltage to the electrodes, gaseous ions can be two dimensionally trapped at the center of the electrodes. By controlling the RF voltage that is applied, it becomes possible to discharge ions with a specific mass or, conversely, trap only those ions with a specific mass.\nA tandem mass spectroscopy (MS/MS) can be conducted using a quadrupole ion trap, as described in the U.S. Pat. No. Re. 34000. In this apparatus, those ions for which no analysis is required are discharged prior to MS/MS. Namely, the removal of the ions for which no analysis is required is not conducted prior to the primary mass spectroscopy. A RF voltage that resonates with the ions is then applied in order to increase the kinetic energy. As a result of these operations, dissociated ions (fragment ions) are created by the collision induced dissociation (CID) with remaining molecules. By subjecting these fragment ions to mass spectroscopy (tandem mass spectroscopy), the mass of the fragment ions can be determined. In this case, it is necessary to initially conduct a mass spectroscopy without involving a CID (primary mass spectroscopy) in order to determine the ions as the object of a tandem mass spectroscopy (MS/MS, or a secondary mass spectroscopy). It is also possible to repeat a similar operation to further conduct a tandem mass spectroscopy (MSn) on a specific dissociated ion.\nRecently, mass spectroscopic methods are often employed for an exhaustive analysis of proteins. Analytical Chemistry, Vol. 73 (2001), pp. 5683-5690, describes examples of analysis called a shotgun analysis. In this technique, a peptide mixture prepared by subjection a protein to enzymatic digestion is separated using a liquid chromatograph, and a separated sample is then subjected to a tandem mass spectroscopy using a quadrupole ion-trap mass spectrometer. With reference to the determined mass of the ion and that of the fragment ion, a database of proteins or genes is searched in order to identify a protein. In case the types of the peptide mixture are too numerous, each peptide might not be completely separated in the liquid chromatograph, and a plurality of kinds of peptides might be simultaneously introduced into the mass spectrometer. This gives rise to the need for automatic tandem analysis called data-dependent analysis. Specifically, the band width of a separated sample separated in a liquid chromatograph is in the order of one minute, and the number of kinds of ions that can be subjected to tandem mass spectrometer at one time is limited to five. In many cases, the ions with greater ion intensities are preferentially subjected to tandem mass spectroscopy, although this depends on the setting of the data-dependent analysis.\nA technical material for the quadrupole ion-trap mass spectrometer manufactured by ThermoFinnigan (www.thermo.com/eThermo/CMA/PDFs/Articles/articlesFile—10918.pdf) describes a dynamic exclusion function. Prior to the start of analysis, the masses of those ions to be excluded from tandem mass analysis are entered and then a list is prepared. By this operation, it becomes possible to exclude those ions put on the list as the objects of data-dependent analysis (tandem mass spectrometer). When this function is to be employed, a conventional mass spectroscopy is conducted first without involving the CID, and then the mass of the ions to be detected is determined. Next, priorities of the ions as the objects of tandem mass spectroscopy are determined in the detected ions, whereupon those ions put on the list are excluded from the objects of data-dependent analysis (tandem mass spectroscopy)."}
{"text": "The Internet has become a popular place to conduct business. Through Web auction sites, Web sites for displaying classified ads, Web shopping malls, online chat rooms, and other online transaction facilitation sites, two consumers may agree to a transaction. Frequently, such transactions involve the exchange of goods or services for money. While consumers frequently find that agreeing to transactions on the Internet is easy, completing a payment to consummate the transaction is more difficult.\nTypically, two consumers who have agreed through the Internet to exchange goods for money resort to offline methods to perform the exchange. For example, the seller may ship the goods to the buyer through a shipping service, and the buyer may send a paper check to the seller.\nSuch offline methods of exchange are problematic. Because the buyer and the seller are usually strangers, they may not trust each other to perform their mutual obligations under the agreement. Accordingly, they may be unable to agree whether the buyer will send the check first or the seller will send the goods first. Even if the buyer and the seller agree that the seller will ship the goods at the same time as the buyer sends the check, the seller has no guarantee that the check will not bounce. Likewise, the buyer has no guarantee that the goods will arrive in satisfactory condition. Accordingly, a significant percentage of transactions to which an individual buyer and seller have agreed upon over the Internet are never consummated.\nAnother inconvenience of transactions agreed upon by individuals over the Internet is that the buyer is often limited to paying by cash or paper check. More convenient payment instruments exist, such as credit cards and bank account debits through electronic fund transactions. However, the buyer typically does not have the option to use these other payment instruments when the seller is an individual as opposed to a retail business that has been pre-established as an online merchant.\nThe term “merchant” is used herein to refer to a seller of goods or services who is authorized by a credit card association (such as DISCOVER, VISA, or MASTERCARD) to submit to the credit card association charges on credit cards belonging to members of the credit card association. After receiving an authorization for the charge, the merchant then receives from the credit card association a direct deposit into the merchant's bank account of the amount of the charge. As known to those skilled in the art, a business must undergo an approval process in order to become a merchant, and upon approval, the merchant is assigned a merchant number.\nAlthough retail businesses are routinely set up as merchants in order to accept payments through credit cards or electronic fund transactions, this is not an adequate solution to facilitating payments between individuals over the Internet. For example, merchants, after undergoing an extensive underwriting effort, are typically given special privileges, such as a general authorization to charge credit cards. This general authorization provides the merchant with the ability to commit fraud. Specifically, the merchant is capable of charging a customer's credit card more than he should. Also, the merchant may submit charges on a credit card belonging to a credit card association member with which the merchant has never had any contact. For these reasons, the idea of allowing individual sellers to become merchants has heretofore been rejected.\nAnother problem with an individual seller becoming a merchant is that the approval process for becoming a merchant is frequently more of a hassle than an occasional seller is willing to undergo. The purpose of the approval process is to reduce the risk of fraud by the merchant. Accordingly, the seller usually must submit extensive background information for consideration in the approval process. This may be inconvenient and time consuming for the seller.\nTherefore, there is a need in the art for a safe and convenient method by which one consumer can pay a second consumer over the Internet."}
{"text": "The broad field of molecular electronics was introduced in the 1970's by Aviram and Ratner. Molecular electronics achieves the ultimate scaling down of electrical circuits by using single molecules as circuit components. Molecular circuits comprising single molecule components can function diversely as switches, rectifiers, actuators and sensors, depending on the nature of the molecule. Of particular interest is the application of such circuits as sensors, where molecular interactions provide a basis for single molecule sensing. In particular, informative current changes could include an increase, or decrease, a pulse, or other time variation in the current.\nNotwithstanding the achievements in the field of molecular electronics, new molecular circuits that can function as molecular sensors are still needed. In particular, the need still exists for improved single molecule systems that can yield molecular information with greater signal-to-noise ratios such that signals truly indicative of molecular interactions are distinguishable from non-informative noise."}
{"text": "In conventional DIAMETER networks, there is currently no way for a DIAMETER node, such as a DIAMETER relay node or DIAMETER routing node, to automatically receive information from a peer DIAMETER node indicating which originating hosts or originating realms the peer DIAMETER node should receive particular types of traffic from. Expressed another way, there is currently no mechanism in the DIAMETER specifications that allows a peer DIAMETER node to specify, to a DIAMETER relay or routing node, that when the DIAMETER relay or routing node receives traffic from a particular originating host or originating realm, the traffic should be directed to the peer DIAMETER node. Internet engineering task force (IETF) request for comments (RFCs) 3588 and 3589 are incorporated herein by reference in their entireties. For example, in the event that network operator X signs an agreement with operator Y that all traffic transmitted through hub provider A's network should be routed to operator X's realm, it may be desirable for operator X to automatically notify hub provider A of the agreement between operator X and operator Y so that hub provider A can update its routing tables accordingly.\nCurrently, in the exemplary scenario described above, hub provider A must manually gather origin-based routing information and update its routing tables. For example, a human user may be required to manually log into a terminal associated with one or more DIAMETER nodes in the hub provider's network and add origin-based rules to the routing tables.\nWhile no method currently exists for automatically notifying and updating the routing tables of a DIAMETER node with origin-based routing information, conventional solutions exist for automatic notification and routing table updating for destination-based routing information. Currently, this is accomplished using the domain name system (DNS). For example, if network operator X wishes to know the networks (i.e., realms or hosts) to which he can route traffic, operator X sends a DNS query to a DNS server that maintains this information. Based on the information included in the DNS response, network operator X (i.e., the DIAMETER node in realm X that initiated the DNS query) may automatically update its routing tables with the destination-based routing information.\nAs may be appreciated from the above discussion, one drawback to conventional methods for populating or updating routing information in DIAMETER nodes is that automatic routing table population is only available for destination-based routing information. As a result, a corollary drawback of conventional methods is therefore that populating and/or updating routing information for DIAMETER nodes is a manual process that may be slow and prone to error.\nAccordingly, in light of these difficulties, a need exists for improved methods, systems, and computer readable media for automatically populating and/or updating routing tables of DIAMETER nodes with origin-based routing information."}
{"text": "A battery pack for an electric car in which a plurality of battery housing parts are defined by disposing a lattice-shaped partition frame on an upper face of a plate-shaped battery frame for the electric car, and a plurality of batteries are mounted in each of the battery housing parts, is known from Patent Document 1."}
{"text": "1. Field of the Invention\nThe present invention relates to a mobile telecommunications technology, and more particularly, to a technique for receiving and transmitting downlink reference signals. Herein, downlink reference signals for data demodulation are efficiently received and transmitted in a single user mode or a multi-user mode.\n2. Discussion of the Related Art\nIn a mobile telecommunications system, a user equipment (UE) may receive information from a base station through a downlink, and the UE may also transmit information through an uplink. Information transmitted or received by the UE may include data and diverse control information. And, depending upon the type and usage of the information transmitted or received by the UE, a variety of physical channels may exist.\nFIG. 1 illustrates a general view showing the physical channels used in a mobile telecommunications system, such as a 3rd generation partnership project (3GPP) long term evolution (LTE) system and a general method for transmitting signals. When power of a UE is turned off and then turned back on, or when using a UE newly introduced to a cell, in step 101, the UE performs an initial cell search process in order to be in synchronization with the base station. In order to do so, the UE receives a primary synchronization channel (P-SCH) and a secondary synchronization channel (S-SCH) from the base station, so as to be in synchronization with the base station, thereby being able to acquire information such as cell ID. Thereafter, the UE receives a physical broadcast channel from the base station, thereby being capable of acquiring broadcast information within the cell. Meanwhile, during an initial cell searching step, the UE receives a downlink reference signal (DL RS), thereby being able to verify the downlink (DL) channel status. After completing the initial cell search, in step 102, the UE may receive a physical downlink control channel (PDCCH) and a physical downlink shared channel (PDSCH) based upon the physical downlink control channel information, so as to acquire more detailed system information.\nMeanwhile, in case the UE has not completed its access to the base station, the UE may perform a random access procedure, as shown in step 103 to step 106, in a later process in order to complete its access to the base station. For this, the UE transmits a characteristic sequence as a preamble through a physical random access channel (PRACH) (S103). Then, the UE may receive a response message respective to the random access through the physical downlink control channel (PDCCH) and its corresponding physical downlink shared channel (PDSCH) (S104). Subsequently, with the exception of a handover, in case of a contention-based random access, the UE may perform a contention resolution procedure, such as transmitting additional physical random access channels (PRACHs) (S105) and receiving the respective physical downlink shared channels (PDSCHs) (S106). After performing the above-described procedure, the UE may receive physical downlink control channel (PDCCH)/physical downlink shared channel (PDSCH) (S107) and may transmit physical uplink shared channel (PUSCH)/physical uplink control channel (PUCCH) (S108), as a general uplink/downlink (UL/DL) signal transmission procedure.\nFIG. 2 illustrates a block view showing a signal processing procedure for transmitting a downlink signal from a base station. In the 3GPP LTE system, the base station may transmit at least one or more code words via downlink. Each of the at least one or more code words may be processed through a scrambling module 301 and a modulation mapper 302 as a complex symbol. Thereafter, the complex symbol is mapped to multiple layers by a layer mapper 303. Herein, a precoding module 304 multiplies each layer by a selected precoding matrix depending upon the channel status, thereby allocating (or assigning) the processed layers to each transmission antenna. Each transmission signal processed as described above for the respective antenna is mapped to a time-frequency resource element, which is to be used by a resource element mapper 305 for transmission. Subsequently, each of the transmission signals passes through an OFDM signal generator 306 so as to be transmitted through the respective antenna.\nHereinafter, a downlink reference signal that is used in the 3GPP LTE system will be described in detail. The 3GPP LTE system uses antenna number 0 to antenna number 5 as its logical antenna ports. Herein, each antenna port is not divided (or classified) by a physical division (or classification). Therefore, the question of mapping each logical antenna index to which actual physical antenna index would relate to the implementation by each manufacturer.\nIn the 3GPP LTE system, three different types of reference signals are used as downlink reference signals. The three types include cell-specific reference signals (non-associated with MBSFN transmission), MBSFN reference signals associated with MBSFN transmission, and UE-specific reference signals. A cell-specific reference signal corresponds to a reference signal generated by using a cell ID for each cell as an initial value. Herein, antenna port 0 to antenna port 3 may be used for transmitting the cell-specific reference signals. An MBSFN reference signal is used for acquiring downlink channel information respective to the MBSFN transmission. Herein, antenna port 4 may be used for transmitting the MBSFN reference signal. Meanwhile, in the 3GPP LTE system, a UE-specific reference signal is supported for a single antenna port transmission of the PDSCH. Herein, antenna port 5 may be used for transmitting the UE-specific reference signal. The UE may receive from an upper layer (or higher layer) (e.g., a MAC layer or higher) information on whether such user-specific reference signals exist so as to be used for PDSCH demodulation.\nFIG. 3 illustrates an example of a specific reference signal being mapped in a time-frequency resource region and transmitted, when the 3GPP LTE system uses a general cyclic prefix. Referring to FIG. 3, the horizontal axis represents a time region, and the vertical axis represents a frequency region. In the time-frequency region shown in FIG. 3, the smallest squared region corresponds to 1 OFDM symbol in the time region and to 1 subcarrier in the frequency region. In the 3GPP LTE system, when a normal cyclic prefix (CP) is used, one slot includes 7 OFDM symbols, and one sub-frame includes 2 slots. FIG. 3 illustrates a pattern where a UE-specific reference signal being transmitted through antenna port 5 is mapped to the time-frequency region throughout even-numbered slots and odd-numbered slots, thereby being transmitted."}
{"text": "Today, integrated circuit's sensors can require a level of precision which cannot be reached with current ultra large scale integration (ULSI) production techniques.\nA conventional solution consists in measuring device deviations with respect to a target value once the wafer manufacturing is complete, and in compensating deviations accordingly via extra processing. Depending on the device type, different solutions are currently in use. For example, digital coding is adopted via fuse concepts with fine-tuning of the electrical characteristics of the product after its final electrical testing at wafer-level. For another example, current sensors are embedded in many ICs to insure constant monitoring and protection of the device during circuit start-up or malfunction. At least some effects of implementing the teaching disclosed herein are as follows: An alternative to conventional laser fuses is provided that enables use of fuses where a conventional laser fuse cannot be used thus opening the processing to a broader range of applications. During wafer processing, techniques that are disclosed herein allow to have chip-selective modifications provided on a wafer, whereby individual chips can be manufactured while using lithography masks designed for production of multiple non-individual chips on a same wafer.\nAt least some effects associated with the introduction of post-processing printed structures can be as follows: Conventional process technologies can be easily adapted or complemented according to some implementations so as to perform digital coding without necessitating Laser tools. Thus, negative effects typically associated with a use of conventional Laser such as alignment problems can be avoided."}
{"text": "The subject matter of the present application relates to microelectronic assemblies and fabrication methods therefor, and more particularly to the structure of and fabrication method for a multilayer interconnect element.\nThere is a current need for microelectronic interconnect elements to provide greater wiring density. Microelectronic interconnect elements include, for example, package substrates used for direct interconnection to microelectronic elements such as semiconductor chips. Other types of interconnect elements include circuit panels which can be directly connected to microelectronic elements or indirectly, such as through a package substrate of a packaged chip. The need is felt especially to improve the density of metal wiring lines, e.g., conductive traces on a dielectric element, as measured by the pitch of the metal lines and minimum spacing between adjacent metal lines.\nSome package substrates and circuit panels have multiple dielectric layers and metal wiring lines provided on some or all of the dielectric layers. A multi-layer wiring substrate 12 according to the prior art, referred to as “High Density Interconnect” is illustrated in FIGS. 1-2. The substrate 12 has a plurality of dielectric layers, two such dielectric layers 14, and 14′ being shown in FIG. 1. As shown therein, each of a plurality of metal lines 10, 10′ and 10″ has approximately the same width w and thickness t.\nOne limitation of the substrate shown in FIG. 1 is a vertical distance factor d by which each of the metal lines 10, 10′ and 10″ is spaced from closest adjacent metal lines (of lines 10, 10′ and 10″) in a vertical direction 30, i.e., the direction of the thickness of each metal line. Each of the metal lines 10 and 10′ is supported by a respective dielectric layer 14 or 14′. As illustrated in FIG. 1, the metal lines 10′ and 10″ are separated in a vertical direction 30 of the substrate 12 by a distance d through a portion of the thickness td of the dielectric layer 14′. A minimum vertical spacing constrains the metal wiring density within the volume occupied by metal lines and dielectric layers 14 of the substrate 12. As further shown in FIG. 2, each of a plurality of traces 10″ adjacent to each other in a horizontal direction 40 has width w and is spaced from the adjacent trace 10″ by a spacing s. Thus, a minimum pitch of the traces 10″, measured between the centers of adjacent traces, is the value of w+s. A minimum spacing s is required for manufacturability of the traces. For example, the traces 10″ of FIG. 2 may be formed subtractively by etching a metal layer. In such case, a constraint in the form of a minimum spacing s is imposed by the resolution of the photolithographic exposure process used to define an etch mask, e.g., a photoresist mask, and the need for the etching process to reliably produce separated traces from a metal layer having a given thickness t. In another example, when the traces 10″ of FIG. 2 are formed in an additive manner by electroplating, a minimum spacing s is imposed by the resolution of the photolithographic exposure process used to define a plating mask, e.g., a photoresist mask, the electroplating process used to form the lines, and the requirements of processes employed after the plating process, e.g., photoresist mask removal. Accordingly, in a HDI implementation, the resulting multi-layer substrate 12 has adjacent traces 10″ spaced apart in a horizontal direction 40 of the substrate by a minimum spacing s. Also, a minimum distance d separates traces of adjacent dielectric layers in a vertical direction 30 of the substrate."}
{"text": "This invention generally relates to overhead doors and, more specifically, to an improved tension spring counterbalancing mechanism use in overhead doors, such as residential garage doors.\nOverhead doors generally require a counterbalancing force which enables the door to be more easily moved between opened and closed positions either manually or by way of a powered opening device. Often, overhead door systems rely on one or more extension springs, placed in tension when the door is in a closed position, so as to provide the desired counterbalancing force. In tension spring counterbalancing systems, the tension in the spring is released as the door is lifted thereby effectively reducing the weight of the door which must be lifted either by the motor of a door opener or by hand. These extension springs must be stretched or extended during the installation of the overhead door such that they are supplied with the necessary counterbalancing tension.\nPresently, the installation of extension spring systems involve labor intensive procedures on the part of the installer. That is, in order for the installer to stretch or extend the spring, the full weight of the overhead door, which may be 200-300 pounds, must be manually lifted and clamped or otherwise propped up in a fully opened position. This procedure usually requires three people, i.e., two people for lifting the door and one more for clamping it in the open position. The operating cable of the door is then attached to the free end of the spring. Since the other end of the spring is fixed, the spring is tensioned when the door is lowered to the closed position. This type of system is not only difficult to install but the door may be inadvertently released from the open position when the spring or springs have not yet been attached to the operating cable. These concerns are amplified by the fact that extension spring systems are often installed by residential homeowners who, in general, are inexperienced with overhead door installation.\nAnother problem associated with overhead doors utilizing extension spring counterbalancing systems concerns the prevention of damage and/or injury resulting from the spring breaking while under tension. Past solutions to this problem have generally involved the use of a rope or cable passed centrally through the spring. Such safety cables are disclosed in U.S. Pat. Nos. 3,958,367; 4,082,133; and 4,640,049. In theory, spring pieces are meant to remain suspended on the rope or cable if the spring should break under tension. However, there is no guarantee that when a spring under tension breaks, it will break in only one place at a location in which the two remaining pieces of the spring will be adequately suspended on the rope or cable. The spring may, for example be brittle and break into several small pieces which would not be retained by a rope or cable.\nU.S. Pat. No. 4,757,853 discloses the combined use of a safety rope extending centrally through the spring and a pair of metal end caps placed over the two ends of the spring. The purpose of the two end caps is stated to be for containment of the extreme ends of the spring through which the safety rope does not extend. While such end caps may provide containment for the ends of the spring, this safety system still does not provide full protection for the entire length of the spring and, in addition, adds undesirable complexity to the overhead door system.\nA need in the art therefore exists for improvements in the art of overhead doors employing extension springs. Specifically, an extension spring system is needed which allows easier installation of the overhead door, especially with regard to the procedures involved with supplying counterbalancing tension to the spring. Furthermore, a system is needed which provides for full containment of the extension spring to guard more completely against damage and injury in the event that the spring breaks under tension."}
{"text": "Pipeline network processors are designed to forward Internet Protocol (IP) packets at an extremely high data rate in excess of two million packets per second. Pipeline network processors typically have a limited number of instruction cycles and memory to perform the applicable task. Implementing features or functions that do not fit the traditional packet forwarding model is a challenging endeavor and may have a negative impact on normal packet forwarding functions such as packet filtering and quality of service processing.\nWith respect to cable modem routers, one of the required and commonly used features is Multicast Echo. This feature enables cable media to behave like a standard shared media, as in Ethernet, in a local area network environment for Internet Protocol multicast traffic. Cable is a unidirectional media wherein packets are sent from customer premises equipment through a cable modem to an upstream port of a cable modem termination system. Packets are also sent from a downstream port of the cable modem termination system through the cable modem to the customer premises equipment. From a cable standpoint, the configured combination of one or more upstream ports and one or more downstream ports is known as a Media Access Control (MAC) Domain which is a cable media equivalent of a local area network segment.\nIn a local area network environment, a device is able to receive all multicast packets sent by other devices on the same segment. However in a cable environment, only half of the multicast processing has been completed when a multicast packet is received at an upstream port. The other customer premises equipment connected to the downstream port of the MAC Domain did not receive the multicast packet. Though other devices on the downstream port may be known to be present, normal multicast processing does not add the input interface to the list of output interfaces for multicast traffic and thus no relationship information among ports. This is intentionally done and required by the multicast protocol to avoid multicast traffic loops. The Multicast Echo feature overrides normal multicast processing and permits multicast packets to be forwarded to downstream ports of the MAC Domain.\nWhen a high speed, high end cable modem termination system is developed using a pipeline network processor to forward packets between a large number, several hundred, of cable interfaces and a large number of trunk interfaces, implementing the Multicast Echo feature for the cable interfaces becomes extremely difficult. The cable line cards strip off the Data Over Cable (DOCSIS) header before passing the packet to the pipeline network processor. The pipeline network processor needs to add the original DOCSIS header to the multicast packet in order to forward the packet to one or more downstream ports. There needs to be support for tens of thousands of multicast groups, different service flows, and many MAC Domains. The amount of memory necessary to hold such large data structures and the processing cycles required to search them would preclude the use of a pipeline network processor in terms of cost and performance. A single multicast flow, which is less than 1% of the total traffic, could severely impact the processing of the other 99% of the traffic. Therefore, it is desirable to implement a Multicast Echo feature using a pipeline network processor despite its limited resources."}
{"text": "A data structure is a particular way of storing and organizing data in a computer so that it can be used efficiently. Data structures are used in almost every program or software system. There are many different types of data structures, such as hash tables, binary search trees, skip lists, ternary search trees, etc. In particular, a hash table is a data structure that uses a hash function to map identifying values, known as keys (e.g., a person's name) to their associated values (e.g., their telephone number). Currently, in many situations, hash tables turn out to be more efficient than search trees or other table lookup structures. For this reason, they are widely used in many kinds of computer software, particularly for associative arrays, database indexing, caches and sets.\nWhile hash tables are currently more efficient than search trees or other table lookup structures, the memory footprint (referring to the amount of main memory that a program uses or references while running) and access speed (referring to the lookup rate) could be improved. Hence, the functionality of hash tables should be maintained while reducing the memory footprint and improving access speed."}
{"text": "Tyrosine ammonia lyase (TAL; also referred to as tyrase, L-tyrosine ammonia lyase, and “L-tyrosine ammonia lyase [trans-p-hyroxycinnamate forming]”), along with histidine ammonia lyase (HAL) and phenylalanine ammonia-lyase (PAL) are members of the aromatic amino acid lyase family (EC 4.3.1.23-1.25 and 4.3.1.3). The enzymes having TAL activity are currently classified in EC4.3.1.23 (previously classified as EC 4.3.1.5). TAL catalyzes the formation of p-coumaric acid from L-tyrosine.\nTyrosinemia (also referred to as “hereditary tyrosinemia,” and “hypertyrosinemia”) is a genetic disorder characterized by elevated blood levels of tyrosine, due to the deficiency of an enzyme required for the catabolism of tyrosine in the liver. If untreated, tyrosine and other metabolites accumulate in the tissues and organs of affected individuals, resulting in serious medical issues. Tyrosinemia is an inborn error of metabolism inherited in an autosomal recessive pattern. There are three types of tyrosinemia, each caused by the deficiency of a different enzyme. Currently used treatment methods depend upon the type of tyrosinemia involved. A low protein diet is often used.\nType I tyrosinemia (also referred to as “FAH deficiency,” “fumaryl acetoacetase deficiency,” “fumaryl aceotacetate hydrolase deficiency,” “hereditary infantile tyrosinemia,” and “hepatorenal tyrosinemia”) is caused by a deficiency of fumarylacetoacetate hydrolase, due to mutations in the fah gene. This is the most severe form of the disease, with symptoms usually appearing in the first few months of life, commonly including failure to thrive, diarrhea, bloody stools, vomiting, jaundice, enlarged liver, the tendency to easily bruise, lethargy, irritability, fever, and other symptoms, such as a distinctive cabbage-like odor of the skin and urine. Some affected infants have repeated neurologic episodes of acute polyneuropathy, characterized by severe leg pain, as well as altered mental status, abdominal pain, and respiratory failure. Infants with the acute form are typically affected at birth and there is a rapid onset of symptoms that can lead to developmental delays, enlarged spleen, ascites, kidney disease, and blood clotting abnormalities. Untreated, it can lead to hepatic and renal failure, nervous system problems, and an increased risk of liver cancer (e.g., hepatocellular carcinoma). In some cases, hypertension and hypertrophic cardiomyopathy are present. If untreated, this disease can be fatal. In the less-common chronic form, the symptoms exhibit a more gradual onset and tend to be less severe. Affected infants initially exhibit vomiting, diarrhea, enlarged liver and spleen, and failure to thrive. Eventually, progressive liver cirrhosis occurs, leading to chronic liver failure, developmental delays, and renal Fanconi syndrome (a rare kidney disorder characterized by weakening and softening of the bones [rickets], vomiting, dehydration, weakness, and fever). In some cases, the most effective treatment has been full or partial liver transplant. Worldwide, this form affects approximately 1 in 100,000 human births (Genetics Home Reference, U.S. National Library of Medicine).\nType II tyrosinemia (also referred to as “keratosis palmoplantaris-corneal dystrophy,” oculocutaneous tyrosinemia,” “Richner-Hanhart syndrome,” “tyrosinemia due to TAT deficiency,” and “tyrosinema due to tyrosine aminotransferase deficiency,”) is caused by a deficiency of tyrosine aminotransferase, due to mutations in the tat gene. It affects the eyes, skin, and mental development. As with Type 1 tyrosinemia, symptoms usually begin in early life, and include excessive tearing, photophobia, eye pain and redness, and painful skin lesions on the palms and soles. About half of affected individuals have some level of intellectual disability. This form occurs in less than 1 in 250,000 persons (Genetics Home Reference, supra).\nType III tyrosinemia (also referred to as “tyrosinemia due to 4-hydroxyphenylpyruvate dioxygenase deficiency,” “tyrosinemia due to 4-hydroxyphenylpyuriv acid oxidase deficiency,” and “tyrosinemia due to HPD deficiency”) is a rare disorder, caused by a deficiency of 4-hydroxyphenylpyruvate dioxygenase, due to mutations in the hpd gene. Symptoms of this form include intellectual disability, seizures, and intermittent ataxia. This form is very rare, only a few cases have been reported (Genetics Home Reference, supra).\nThere are additional cases in which there are temporary elevated tyrosine levels, due to non-genetic factors such as vitamin C deficiency or premature birth, which results in immature liver enzymes. Differential diagnoses are used to differentiate these transient cases from tyrosinema I, II, or III.\nIn addition to tyrosinema, there are other diseases associated with insufficient or absent tyrosine metabolism. For example, alkaptonuria also referred to as alcaptonuria, is a disease caused by deficiency of homogentisate 1,2-dioxygenase, which is an enzyme involved in tyrosine degradation. This enzyme is encoded by the HGD gene. Insufficient activity of this enzyme results in the accumulation of homogentisic acid. Excess homogentisic acid and related compounds are deposited in connective tissues, causing the cartilage and skin to darken. Over time, arthritis may result due to the accumulation of homogentisic acid and related metabolites in the joints of affected individuals. Homogentisic acid is also excreted in urine, making the urine turn black. Alkaptonuria is a rare disease that affects 1 in 250,000 to 1,000,000 people worldwide (See, Genetics Home Reference, supra).\nTreatment of these diseases has largely been the life-long use of a methionine-, phenylalanine-, and tyrosine-restricted diet. Treatment with nitisinone (NTBC; 2-(2-nitro-4-trifluoromethylbenzol)-1,3-cyclohexane dione; Orfadin®) has been reported to be helpful for type I tyrosinemia and alkaptonuria, due to its inhibition of the 4-hydroxyphenylpyruvate oxidase pathway. However, NTBC must be used in combination with a challenging and costly methionine-, phenylalanine-, and tyrosine-restricted diet to prevent both liver failure and carcinogenesis. There remains a need in the art for easy to administer, effective treatment(s) to ameliorate the symptoms of these diseases and allow patients to utilize normal diets."}
{"text": "1. Field of the Invention\nThis invention relates to improvements in the encoding scheme of gray pixels. More particularly, this invention relates to gray pixel encoding schemes which are used in combination with a halftoning imaging that produces gray level halftoned images.\n2. Description of Related Art\nXerographic printers use halftoning processes to generate digital images from originals. The image is formed by using pixels which are completely black, completely white or different shades of gray. These are represented as explicit gray pixels. When viewed from a distance, the pixels blend together to form the image.\nIn traditional binary xerographic printers, halftone methods use binary systems such that the laser has only two laser intensity levels: ON (black) and OFF (white). Another method is disclosed in U.S. Pat. No. 4,868,587 (Loce et al.), which uses two gray levels: the black and the white level for each pixel. The different levels of gray are based on the intensity of the laser light. Therefore, black has a higher light intensity than either of the gray levels. This method is an improvement over using only black and white pixels. The disadvantage of the system is that the control of the graying levels is difficult and the information is not easily transferable to other printers.\nThe quality of a binary xerographic printer is based on two important features: the halftone frequency, which is the number of halftone cells per linear inch; and the number of distinguishable gray steps. A typical number of halftone cells for good quality image is between 100 to 200 cells per inch. Usually around nine pixels are used per halftone cell. The maximum number of gray steps is limited to the number of pixels per halftone cell plus one. Therefore, a 3.times.3 halftone cell has about 10 output gray steps.\nIn the prior art, a halftone cell was divided into pixels which were turned either on or off to form reflectance modulation level. From a distance, the human eye detects a reflectance of about 0.5%, which would be interpreted as a gray level. In a good quality printer, the number of distinguishable gray steps should be around 100.\nHalftone images are commonly used in printed materials to reproduce continuous tones using printers that are binary in nature. Conventional digital halftoning processes grow halftone dots by activating individual pixels within a halftone cell. The pixels are typically at printer resolution or at some small integral subdivision thereof. The printer resolution is broken down to the number of spots which are used. For example, a 400 SPI printer has a spot size of 1/400th of an inch. Pixels are used as data within the spot boundary."}
{"text": "Most all types of commercially-available power diodes having high reverse breakdown voltage capabilities have N-type bottomside cathodes. A rare exception is the so-called “inverse diode” or “reverse diode” that is commercially available from IXYS Corporation, 1590 Buckeye Drive, Milpitas, Calif. These unusual diodes have P type isolation structures involving a bottomside P type anode region as well as P type peripheral sidewall diffusion regions. Not only do these diodes have very high reverse breakdown voltages, but they also typically exhibit superior dynamic robustness. An attempt was made to extend this “inverse diode” technology to so-called “fast diodes” having lower reverse recovery times. The reverse recovery time of a diode is denoted Trr in the literature and in data sheets. As set forth in U.S. Pat. No. 8,716,745, an N− type epitaxial silicon layer was grown on a P type wafer. The resulting inverse diode was simulated to have superior stability and a high reverse breakdown voltage while at the same time having a thinner N− type layer as compared to a conventional diode having the same reverse breakdown capabilities. For additional information on inverse diode structures and on P type isolation structures, see: 1) U.S. Pat. No. 7,442,630, entitled “Method For Fabricating Forward And Reverse Blocking Devices”, filed Aug. 30, 2005, by Kelberlau et al.; 2) U.S. Pat. No. 5,698,454, entitled “Method Of Making A Reverse Blocking IGBT”, filed Jul. 31, 1995, by N. Zommer; 3) J. Lutz et al., “Semiconductor Power Devices”, pages 146-147, published by Springer, Berlin and Heidelberg (2011); 4) the data sheet entitled “Diode Chip”, DWN 17-18, by IXYS Corporation of Milpitas, Calif. 95035, USA; 5) U.S. Pat. No. 9,590,033, entitled “Trench Separation Diffusion For High Voltage Device”, filed Nov. 20, 2005, by Wisotzki et al.; 6) U.S. Pat. No. 4,351,677, entitled “Method of Manufacturing Semiconductor Device Having Aluminum Diffused Semiconductor Substrate”, filed Jul. 10, 1980, by Mochizuki et al.; 7) U.S. Pat. No. 6,507,050, entitled Thyristors Having A Novel Arrangement of Concentric Perimeter Zones”, filed Aug. 16, 2000, by Green; 8) U.S. Pat. No. 6,936,908, entitled “Forward and Reverse Blocking Devices”, filed Mar. 13, 2002, by Kelberlau et al.; 9) U.S. Pat. No. 7,030,426, entitled “Power Semiconductor Component in the Planar Technique”, filed Mar. 14, 2005, by Neidig; 10) U.S. Pat. No. 8,093,652, entitled “Breakdown Voltage For Power Devices”, filed Aug. 27, 2003, by Veeramma et al.; 11) the 2004 description entitled “FRED, Rectifier Diode and Thyristor Chips in Planar Design”, by IXYS Semiconductor GmbH, Edisonstrasse 15, D-68623, Lampertheim, Germany."}
{"text": "1. Field of the Invention\nThis invention relates to a decoder circuit, for example for use in memory addressing.\n2. Background Art\nFIG. 1 shows an example of the structure of a static random access memory (SRAM) and its addressing circuitry. The memory 1 comprises a grid of memory cells 2 arranged in rows and columns. In this example there are four rows and four columns but typical SRAMs have many more rows and columns. All the cells in a row are connected to a single wordline 3. All the cells in a column are connected to a single bitline 4. Each cell has two stable states, representing high and low output values. When one of the wordlines is activated each bitline takes on the state of the cell to which that bitline is connected and which is also connected to the activated wordline. The bitlines are connected to a multiplexer 5 which allows a single one of the bitlines to be selected to provide the final output from the memory at 6.\nEach cell in the memory is allotted a number. To access a cell in the memory the cell's number is applied in binary format to input lines 7. A decoder 8 takes some of the lines 7 (indicated at 9) as input and thereby determines which one of the wordlines should be activated to access the cell. The remainder of the lines 7 (indicated at 9) pass to the multiplexer 5, which thereby determines which of the bitlines should be selected so that the final output at 6 takes on the state of the desired cell.\nFIG. 2 shows one example of a structure for decoder 8. The lines 9 each branch through an inverter 10 so that address signals representing inverted (indicated at 11) and non-inverted (indicated at 12) versions of each of the lines 9 are available. The appropriate inverted or non-inverted version of each of the lines 9 is applied to a NAND gate 13 corresponding to each wordline. The output of each NAND gate passes to an inverter 14 which drives the respective wordline.\nBy means of the correct set of connections to each of the NAND gates a wordline is activated only when the appropriate set of inputs is applied at 9. This circuit is simple, but has a number of disadvantages if the number of address line inputs is increased.\n1. In many important applications—for example in the caches of high-speed processors—speed of access to the desired memory cell is crucial. However, the large NAND gates are relatively slow.\n2. In almost all applications—and especially in battery-powered applications power consumption is very important. However, in the system of FIG. 2 as the input lines 9 change values potentially many, or all, of the distributed address signals will change, giving rise to a high current consumption. The current consumption is further increased by the need to generate an inverted version of each input line to the decoder.\nIn an alternative structure for decoder 8 a pre-decode level is-added so that fewer distributed signals can change as the input lines 9 change values. In the predecode level the input lines 9 are split into groups that are processed by a first level of NAND gates. The outputs from those NAND gates pass to a second level of NAND gates whose outputs are inverted to drive the wordlines. Since only one of the lines that connect the two levels of NAND gates changes for each change in input values the maximum power consumption is less than for the structure of FIG. 2. Two levels of processing in the alternative structure make it relatively slow for small numbers of addresses, but for wider decoders (e.g. five or more addresses) it can be faster.\nAn alternative solution is to use a precharge decode structure. FIG. 3 illustrates precharge decode circuitry that represents an alternative structure for decoder 8. As before, inverted 11 and non-inverted 12 versions of each of the input lines 9 to the decoder are generated. A decode line 20 is provided for each wordline. Each wordline is driven by an inverter 14 which receives the output from a two-input NAND gate 18. One of the inputs to the NAND gate 18 is a common timing enable signal at 19. The other of the inputs is from the respective decode line 20.\nEach decode line can be taken high by a respective PMOS precharge transistor 21. The precharge transistors are connected with their sources to a high voltage (Voc), their drains to the respective decode line and their gates to a common precharge signal at 22. Each decode line can be taken low by any of a number of NMOS addressing transistors. The addressing transistors 23 are connected with their drains to the respective decode line, their sources to a low voltage (ground) and their gates to a selected one of the inverted and non-inverted input lines 11, 12. The inputs to the gates of the addressing transistors are arranged so that one addressing transistor of each decode line receives a selected inverted or non-inverted version of each of the input lines 9.\nIn use, to decode a set of signals applied to lines 9, a pulse is applied to the precharge line 22 so that the decode lines float at VDD). Then the signals from the input lines 9 are applied to the addressing transistors 23. At all the decode lines corresponding to all the undesired wordlines at least one of the addressing transistors is turned on, so that those decode lines can all discharge to ground. At the decode line corresponding to the desired wordline none of the addressing transistors is turned on, so that decode line continues to float at VDD. After a delay that is sufficiently long for the voltages on all the undesired decode lines to have fallen, a signal is applied to the timing enable line 19, This causes the NAND gate 18 corresponding to the decode line that remains high to produce a low output, whereas all the other NAND gates produce high outputs. By virtue of the inverters 14 this causes the desired one of the wordlines to be selected.\nThe delay before application of the timing enable signal is crucial to the operation of the precharge circuit. The timing enable signal cannot be applied before all the undesired decode lines have discharged to below the input threshold voltage of the NAND gates 18, otherwise more than one wordline will be selected. Therefore, the delay is dependant on the speed with which the decode lines discharge. A decode line with all of its addressing transistors turned on will discharge quickly but a decode line with only one of its addressing transistors turned on will discharge relatively slowly. The timing enable signal is usually derived from the arrangement shown at 24. Another instance of the decode structure is provided, this time with all but one of its addressing inputs connected to ground. The decode line 25 of the arrangement 24 therefore falls towards ground as slowly as any of the undesired decode lines can. The enable input to the NAND gate 26 of the arrangement 24 is held high. As soon as the decode line has fallen sufficiently the output 27 of the NAND gate goes high. The output 27 provides the timing enable signal at 19.\nThe precharge structure of FIG. 3 can be faster than the structures of FIG. 2. However, the discharge of all the undesired decode lines gives rise to high power consumption, and the necessary delay before the timing enable signal is generated reduces the speed of operation of the circuit. The choice between a static decoder (e.g. as shown in FIG. 2) and a precharge decoder is not clear-cut and very often both options have to be investigated when undertaking a design.\nIt would be desirable to have a decoder circuit that allowed for faster decoding, preferably at reduced power consumption. In addition to providing a technically superior solution such a circuit could save a considerable amount of investigatory design work."}
{"text": "The teachings of all the references cited in the present specification are incorporated in their entirety by reference.\nObesity and its associated disorders are common and very serious public health problems in the United States and throughout the world. Upper body obesity is the strongest risk factor known for type-2 diabetes mellitus, and is a strong risk factor for cardiovascular disease. Obesity is a recognized risk factor for hypertension, arteriosclerosis, congestive heart failure, stroke, gallbladder disease, osteoarthritis, sleep apnea, reproductive disorders such as polycystic ovarian syndrome, cancers of the breast, prostate, and colon, and increased incidence of complications of general anesthesia. It reduces life-span and carries a serious risk of co-morbidities above, as well disorders such as infections, varicose veins, acanthosis nigricans, eczema, exercise intolerance, insulin resistance, hypertension hypercholesterolemia, cholelithiasis, orthopedic injury, and thromboembolic disease. Obesity is also a risk factor for the group of conditions called insulin resistance syndrome, or “Syndrome X.”\nIt has been shown that certain peptides that bind to the Y2 receptor when administered peripherally to a mammal induce weight loss. The Y2 receptor-binding peptides are neuropeptides that bind to the Y2 receptor. Neuropeptides are small peptides originating from large precursor proteins synthesized by peptidergic neurons and endocrine/paracrine cells. Often the precursors contain multiple biologically active peptides. There is great diversity of neuropeptides in the brain caused by alternative splicing of primary gene transcripts and differential precursor processing. The neuropeptide receptors serve to discriminate between ligands and to activate the appropriate signals. These Y2 receptor-binding peptides belong to a family of peptides including peptide YY (PYY), neuropeptide Y (NPY) and pancreatic peptide (PP).\nNPY is a 36-amino acid peptide and is the most abundant neuropeptide to be identified in mammalian brain. NPY is an important regulator in both the central and peripheral nervous systems and influences a diverse range of physiological parameters, including effects on psychomotor activity, food intake, central endocrine secretion, and vasoactivity in the cardiovascular system. High concentrations of NPY are found in the sympathetic nerves supplying the coronary, cerebral, and renal vasculature and have contributed to vasoconstriction. NPY binding sites have been identified in a variety of tissues, including spleen, intestinal membranes, brain, aortic smooth muscle, kidney, testis, and placenta.\nNeuropeptide Y (NPY) receptor pharmacology is currently defined by structure activity relationships within the pancreatic polypeptide family. This family includes NPY, which is synthesized primarily in neurons; PYY, which is synthesized primarily by endocrine cells in the gut; and PP, which is synthesized primarily by endocrine cells in the pancreas. These approximately 36 amino acid peptides have a compact helical structure involving a “PP-fold” in the middle of the peptide. Specific features include a polyproline helix in residues 1 through 8, a β-turn in residues 9 through 14, an α-helix in residues 15 through 30, an outward-projecting C-terminus in residues 30 through 36, and a carboxyl terminal amide, which appears to be critical for biological activity. The peptides have been used to define at least five receptor subtypes known as Y1, Y2, Y3, Y4 and Y5. Y1 receptor recognition by NPY involves both N- and C-terminal regions of the peptide; exchange of Gln34 with Pro34 is fairly well tolerated. Y2 receptor recognition by NPY depends primarily upon the four C-terminal residues of the peptide (Arg33-Gln34-Arg35-Tyr36-NH2) preceded by an amphipathic an α-helix; exchange of Gln34 with Pro34 is not well tolerated. One of the key pharmacological features which distinguish Y1 and Y2 is the fact that the Y2 receptor (and not the Y1 receptor) has high affinity for the NPY peptide carboxyl-terminal fragment NPY-(13–36) and the PYY fragment PYY(22–36).\nIt has been shown that a 36 amino acid peptide called Peptide YY(1–36)[PYY(1–36)] [YPIKPEAPGEDASPEELNRYYASLRHYLNLVTRQRY, SEQ ID NO.: 1]. when administered peripherally by injection to an individual produces weight loss and thus can be used as a drug to treat obesity and related diseases, Morley, J. Neuropsychobiology 21:22–30 (1989). It was later found that to produce this effect PYY bound to a Y2 receptor, and the binding of a Y2 agonist to the Y2 receptor caused a decrease in the ingestion of carbohydrate, protein and meal size, Leibowitz, S. F. et al. Peptides, 12:1251–1260 (1991). An alternate molecular form of PYY is PYY(3–36) IKPEAPGEDASPEELNRYYASLRHYLNLVTRQRY [SEQ ID NO.: 2], Eberlein, Eysselein et al. Peptides 10: 797–803, 1989). This fragment constitutes approximately 40% of total PYY-like immunoreactivity in human and canine intestinal extracts and about 36% of total plasma PYY immunoreactivity in a fasting state to slightly over 50% following a meal. It is apparently a dipeptidyl peptidase-IV (DPP4) cleavage product of PYY. PYY3–36 is reportedly a selective ligand at the Y2 and Y5 receptors, which appear pharmacologically unique in preferring N-terminally truncated (i.e. C-terminal fragments of) NPY analogs. It has also been shown that a PYY fragment having only residues 22–36 will still bind to the Y2 receptor. However, if any of the carboxyl terminus of the peptide is cleaved, the peptide looses its ability to bind to the Y2 receptor. Hence a PYY agonist is a peptide, which has a partial sequence of full-length PYY and is able to bind to a Y2 receptor in the arcuate nucleus of the hypothalamus. Hereinafter the term PYY refers to full-length PYY and any fragment of PYY that binds to a Y2 receptor.\nIt is known that PYY and PYY3–36 can be administered by intravenous infusion or injection to treat life-threatening hypotension as encountered in shock, especially that caused by endotoxins (U.S. Pat. No. 4,839,343), to inhibit proliferation of pancreatic tumors in mammals by perfusion, parenteral, intravenous, or subcutaneous administration, and by implantation (U.S. Pat. No. 5,574,010) and to treat obesity (Morley, J. Neuropsychobiology 21:22–30 (1989) and U.S. Patent Application No. 20020141985). It is also claimed that PYY can be administered by parenteral, oral, nasal, rectal and topical routes to domesticated animals or humans in an amount effective to increase weight gain of said subject by enhancing gastrointestinal absorption of a sodium-dependent cotransported nutrient (U.S. Pat. No. 5,912,227). However, for the treatment of obesity and related diseases, including diabetes, the mode of administration has been limited to intravenous IV infusion with no effective formulations optimized for alternative administration of PYY3–36. None of these prior art teachings provide formulations that contain PYY or PYY(3–36) combined with excipients designed to enhance mucosal (i.e., nasal, buccal, oral) delivery nor do they teach the value of endotoxin-free Y2-receptor binding peptide formulations for non-infused administration. Thus, there is a need to develop formulations and methods for administering PYY3–36."}
{"text": "The present invention relates to a tool with a fastener engaging member, and in particular, to a fastener engaging member that is adapted to form an interface with at least one surface on the fastener such that the fastener can be releasably retained to the driving portion of the tool.\nThe prior art has long sought to develop a satisfactory holding attachment for tools that assist the user in holding, piloting and starting a fastener, as well as with the removal of the fastener. One approach is to magnetize the tool. A magnetized tool is only suitable for retaining ferrous fasteners. Magnetized tools also collect ferrous debris, such as metal shavings and chips.\nU.S. Pat. No. 6,302,001 (Karle) discloses a hex-shaped tool head with a circumferential recess to receive as spring washer. The spring washer secures the hex-shaped tool head to the internal surfaces of the screw head. The circumferential recess weakens the tool head.\nU.S. Pat. No. 1,698,521 (Wood); U.S. Pat. No. 1,712,196 (Burger et al.); and U.S. Pat. No. 3,245,446 (Morifuji) disclose a pair of inwardly biased members that grasp the head of the fastener. These devices can typically be used only on fastener with heads within a certain size range. If the fastener head is larger or smaller than that certain size range, the device does not operate as intended. For some of these devices, the shape of the head is also critical to proper operation.\nU.S. Pat. No. 4,016,913 (Anderson) discloses a pair of springs extending between a pair of arms attached to the tool that are adapted to grip the shank or threaded portion of the fastener. The usefulness of the device of Anderson is also limited by the size of the fastener. For large diameter fasteners, longer springs are required. The longer springs, however, are less effective at holding smaller diameter fasteners. Consequently, multiple devices are required to accommodate fasteners with largely varying diameters.\nU.S. Pat. No. 4,197,886 (MacDonald) discloses a fastener holding nosepiece for a driving tool. The nosepiece is removable from the adapter by a quick disconnect feature that permits different nosepieces to be substituted to accommodate fasteners having heads of larger or smaller diameters.\nThe present invention is directed to a tool adapted to releasably retain a fastener. The tool includes a driving portion comprising a plurality of tool surfaces adapted to form an interface with a fastener. At least one polymeric fastener engaging member is located in a recess in the driving portion that extends above one or more of the tool surfaces. The fastener engaging member is adapted to form an interface with at least one surface on the fastener such that the fastener can be releasably retained to the driving portion.\nIn one embodiment, the recess and the fastener engaging member are located in a center region of the tool surface. The size of the center region can vary and may have a surface area larger than the recess and fastener engaging member. In one embodiment, the center region comprises about the middle 70% between the transition edges of adjacent tool surfaces, and more preferably about the middle 50% between the transition edges of adjacent tool surfaces, and most preferably about the middle 30% between the transition edges of adjacent tool surfaces.\nThe present invention is also directed to a driving portion comprising a plurality of tool surfaces adapted to be positioned in the tool receiving recess in a fastener. The fastener engaging member is adapted to form an interface with at least one surface in the tool receiving recess in the fastener such that the fastener can be releasably retained to the driving portion. The present invention is also directed to a tool with a fastener engaging member that is adapted to be positioned around a portion of the fastener.\nIn one embodiment, a single fastener engaging member is attached to the driving portion at only one of the tool surfaces. In another embodiment, a single fastener engaging member is attached to the driving portion along an edge between two adjacent tool surfaces. The fastener engaging member may also extend along the distal end of the tool.\nThe fastener engaging member is located in a recess formed in the driving portion. The recess can be located in one of the tool surfaces or along an edge between two adjacent tool surfaces. Discrete recesses can be located on a plurality of the tool surfaces. In one embodiment, the recess extends through the driving portion such that the fastener engaging member is located in the recess and extends above two non-adjacent tool surfaces on the driving portion. The two non-adjacent tool surfaces are preferably opposing surfaces such that the compressive forces on the fastener engaging member are generally opposing and co-linear.\nIn one embodiment, a reinforcing member is located in the polymeric material. The reinforcing member can be a resilient member that deforms elastically, such as spring member or a wire, or a substantially rigid member. The reinforcing member typically extends above one or more of the tool surfaces of the driving portion. In one embodiment, the reinforcing member extends above the polymeric material. The reinforcing member can also be rigid. In this embodiment, the rigid reinforcing member would be displaced (typically rotated) during compression of the polymeric material.\nThe polymeric material is selected from a group comprising nylon, polypropylene, PVC, ABS, cellulose, acetyl, polyethylene, fluoropolymers, polycarbonate, natural or synthetic rubber, and the like. In one embodiment, the polymeric material comprises an adhesive. The polymeric material typically extends above the tool surface about 0.001 inches to about 0.2 inches, although this distance will vary considerably with the application, such as the type of tool, the type of fastener, the material from which the fastener is constructed, and the like. The tool can be one of a ballpoint tool, a torx(copyright) driver, square drivers, a hex wrench, socket wrench, a flat-head screw driver, a phillips screw driver, an open-ended wrench, a box wrench, or any other tool adapted to releasably engage with a fastener.\nThe present invention is also directed to a tool adapted for use with a fastener having a tool receiving recess. The tool includes a driving portion comprising a plurality of tool surfaces adapted to be positioned in the tool receiving recess. At least one elongated fastener engaging member is located in the recess in the driving portion and extends above one or more of the tool surfaces. The fastener engaging member forms an interface with at least one surface in the tool receiving recess such that the fastener is releasably retained to the driving portion.\nThe fastener engaging member can be a polymeric material, metal, ceramic, or a combination thereof. The fastener engaging member can be configured as a coil spring, a wire, a ribbon, and the like. The fastener engaging member preferably comprises a spring member shaped to generate a biasing force against inside surfaces of the recess where the biasing force retains the elongated fastener engaging member in the recess. A polymeric material, such as an adhesive, can optionally be deposited in the recess with the elongated fastener engaging member.\nThe present method is also directed to a method of forming a tool adapted to releasably retain a fastener. The method includes forming one or more recesses in one or more tool surfaces of a driving portion of the tool. At least one polymeric fastener engaging member is located in each recess such that the fastener engaging member extends above one or more of the tool surfaces.\nThe fastener engaging member can be a polymeric material molded or inserted in the recess. In one embodiment, the driving portion engages with a tool receiving recess on the fastener."}
{"text": "As a bearer network technology for the Next Generation Network (NGN), the Ethernet technology is characterized by low cost, easy operation, and convenient upgrade. The Ethernet technology develops from single networks to a hierarchical and connection-oriented trend, and the current hierarchical Ethernet technology shapes up gradually. In the development process, a Provider Backbone Bridge Traffic Engineering (PBB TE) technology, namely, a Provider Backbone Transport (PBT) technology, is generated.\nThe PBB TE technology implements the connection-oriented feature of the Ethernet, and is a derivative technology of the Ethernet standard. It disables the spanning tree and disables the flood mechanism and the broadcast mechanism of traditional Ethernet, and provides various services by setting up tunnels on the backbone network. The PBB TE technology is developed on the basis of the Media Access Control (MAC)-in-MAC technology. The MAC-in-MAC technology uses the MAC address of an operator to encapsulate a user's MAC address, thus reducing and isolating user MAC addresses in the data transmission process. The MAC-in-MAC technology enables a network to be hierarchical, for example, divides a network into three layers, namely, a user network layer, a provider network layer, and a backbone network layer connected to each provider network. FIG. 1 schematically shows MAC-in-MAC network connections in the prior art. The network connections include a user network layer, a provider network layer, and a backbone network layer. Specifically, the user network layer is composed of a Customer Bridged Network (CBN) X and a CBN Y; the backbone network layer is a Provider Backbone Bridged Network (PBBN) layer; and the provider network layer is composed of Provider Bridged Network (PBN) X and PBN Y, where the PBN X connects the CBN X and the PBBN, and the PBN Y connects the CBN Y and the PBBN.\nThe PBB TE technology sets up a tunnel on the PBBN network to forward the two-layer MAC data frames. The tunnel is identified by a Backbone Destination Address (B-DA) and a Backbone Virtual Local Area Network ID (B-VID) in a B-tag. In the PBB TE tunnel, only the B-tag needs to be identified, and the user information is transparent. The PBB TE technology provides connection-oriented services for data frames, and sets up a working tunnel and a standby tunnel through an outer tag (B-tag)+B-DA (“+” means “plus”). The number of the standby tunnels to be set up depends on the actual network configuration. Each B-VID identifies the working tunnel and the standby tunnel between a Backbone Source Address (B-SA) and a Backbone Destination Address (B-DA). FIG. 2 schematically shows PBB TE tunnel connections in the prior art. The data frames are forwarded between the B-SA and the B-DA1 through the working tunnel identified by B-VID1, and the data frames are forwarded between the B-SA and the B-DA2 through the working tunnel identified by B-VID1, and a standby tunnel is configured for the working tunnel and is identified by B-VID2. The standby tunnel provides end-to-end protection for the working tunnel between the B-SA and the B-DA2. The data frames are forwarded between the B-SA and the B-DA3 through the working channel identified by B-VID1. On a PBBN network, the working tunnel between the B-SA and the B-DA1 is identified by B-DA1+B-VID1; the working tunnel between the B-SA and the B-DA2 is identified by B-DA2+B-VID1, and the standby tunnel is identified by B-DA2+B-VID2; and the tunnel between the B-SA and the B-DA3 is identified by B-DA3+B-VID1. In the PBB TE tunnel, the B-VID is reusable for different B-DAs, but it is necessary to ensure the combination of B-DA and B-VID, that is, the B-DA+B-VID to be unique throughout the PBBN. The PBB TE technology provides a good foundation for network expansion.\nNormally, the data traffic (namely, data frames) sent from the B-SA runs along the working tunnel to the B-DA. When the working tunnel fails, the preconfigured standby tunnel becomes active and takes over the data traffic sent from the B-SA. The data traffic is sent to the B-DA through the standby tunnel to implement protection switching. The protection switching provides carrier grade end-to-end protection between the B-SA and the B-DA. FIG. 3 schematically shows PBB TE end-to-end protection in a prior art. When the working tunnel identified by B-VID1 between the B-SA and the B-DA2 fails, the B-SA switches the data frames over to the preconfigured standby tunnel identified by B-VID2. On the PBBN, the standby tunnel is identified by B-DA2+B-VID2, and the data frames are sent to the B-DA2 through the standby tunnel, thus implementing protection switching and PBB TE end-to-end protection."}
{"text": "1. Field of the Invention\nThe present invention relates to systems and method for disinfecting a surface, such as a keyboard or display screen.\n2. Background of the Related Art\nThe ordinary use of computer keyboards and touchscreens can lead to an accumulation of pathogens growing on these surfaces. These pathogens may then increase in number or be spread between individuals that share use of the keyboards and touchscreens. A “pathogen” may be any infectious agent that can cause disease, such as a virus, bacterium, prion, fungus, viroid or parasite.\nAlthough it is desirable to periodically clean the surfaces of a keyboard and touchscreen, doing so can be difficult and tedious, while posing a threat of damage to the devices. In particular, a keyboard includes a large number of individual surfaces and mechanical connections that can become displaced or damaged due to physical forces, and electrical components that can be damaged by corrosive fluids or excessive amounts of fluids.\nCurrent disinfection methods involve either wiping the keyboard with a cloth carrying some type of disinfectant chemical or manually moving a UV light-emitting wand over the keyboard. The wiping method leaves many areas untouched and some users are uncomfortable wiping their keyboards with a liquid (or gel). The UV light method is rarely used and requires purchase and care of a separate device. Still, neither of these cleaning methods will prevent the spread of pathogens if a user will not take the time to use them properly."}
{"text": "In recent years, optical coherence tomography (OCT) has been drawing attention. The OCT creates an image representing the exterior or interior structure of an object to be measured using light beams from a laser light source or the like. Unlike X-ray computed tomography (CT), the OCT is not invasive on the human body, and therefore is expected to be applied to the medical field and the biological field, in particular. For example, in the opthalmological field, apparatuses for forming images of the fundus oculi or the cornea have been in practical use.\nPatent Document 1 discloses a device using Fourier-domain OCT or frequency-domain OCT. This device irradiates an object to be measured with a beam of low-coherence light, and superimposes the light reflected from the object on reference light to generate interference light. The device then obtains the spectral intensity distribution of the interference light, and applies Fourier transform thereto to acquire an image of the morphology of the object to be measured in the depth direction (z direction). The device includes a galvanometer mirror configured to scan a light beam (measurement light) in a direction (x direction) perpendicular to the z direction, thereby forming an image of a desired area of the object to be measured. The image formed by the device is a two-dimensional cross-sectional image in the depth direction (z direction), taken along the scanning direction (x direction) of the light beam. Such technique using a spectrometer is called “spectral-domain”.\nPatent Document 2 discloses a technology, in which measurement light is scanned in the horizontal direction (x direction) and the vertical direction (y-direction) to thereby form a plurality of two-dimensional cross-sectional images in the horizontal direction. Based on the cross-sectional images, three-dimensional cross-section information is obtained for a measurement range. As the three-dimensional imaging, for example, there are a method of arranging a plurality of cross-sectional images in the vertical direction (referred to as “stack data”, etc.), a method of performing rendering on volume data (voxel data) based on the stack data to thereby form a three-dimensional image, and the like.\nPatent Documents 3 and 4 disclose OCT devices of other types. Patent Document 3 discloses an OCT device, which scans (sweeps) the wavelengths of light irradiated to the object to be measured, and sequentially detects interference light obtained by superimposing reflected light of each wavelength on reference light to acquire spectral intensity distribution. The device applies Fourier transform to the spectral intensity distribution to form an image of the morphology of the object to be measured. Such an OCT device is called swept-source OCT. The swept-source OCT is a type of Fourier-domain OCT.\nPatent Document 4 discloses an OCT device, which irradiates light beams having a predetermined diameter to an object to be measured, and analyzes the components of interference light obtained by superimposing the reflected light on reference light. Thereby, the device captures an image of the object to be measured in a cross-section perpendicular to the traveling direction of the light. Such an OCT device is called full-field OCT or en-face OCT.\nPatent Document 5 discloses a configuration in which OCT is applied to the ophthalmologic field. Incidentally, before the application of OCT, a fundus camera, a slit lamp, a scanning laser opthalmoscope (SLO), or the like has been used as a device for observing the subject's eye (see, for example, Patent Documents 6, 7, and 8). The fundus camera is a device that irradiates the subject's eye with illumination light and receives the light reflected from the fundus to thereby capture an image of the fundus. The slit lamp is a device that cuts out an optical section of the cornea using a slit light to thereby acquire an image of the cross-section of the cornea. The SLO is a device that scans the fundus with a laser beam, and detects its reflected light with a high-sensitivity element such as a photomultiplier tube for imaging the morphology of the fundus surface.\nThe devices using OCT offer advantages with respect to the fundus camera in that they can acquire high-resolution images, and also that they can obtain cross-sectional images as well as three-dimensional images.\nAs described above, the devices using OCT can be used for the observation of different parts of the eye, and are capable of acquiring high-resolution images. Therefore, the OCT devices have been applied to a variety of ophthalmic diseases. For example, there has been a known device, which is made of a combination of an OCT device and a subjective visual acuity test device, and provides materials for the diagnosis of maculopathy and glaucoma (see Patent Document 9).\n[Patent Document 1] Japanese Unexamined Patent Application Publication No. Hei 11-325849\n[Patent Document 2] Japanese Unexamined Patent Application Publication No. 2002-139421\n[Patent Document 3] Japanese Unexamined Patent Application Publication No. 2007-24677\n[Patent Document 4] Japanese Unexamined Patent Application Publication No. 2006-153838\n[Patent Document 5] Japanese Unexamined Patent Application Publication No. 2008-73099\n[Patent Document 6] Japanese Unexamined Patent Application Publication No. Hei 9-276232\n[Patent Document 7] Japanese Unexamined Patent Application Publication No. 2008-259544\n[Patent Document 8] Japanese Unexamined Patent Application Publication No. 2009-11381\n[Patent Document 9] Japanese Unexamined Patent Application Publication (Translation of PCT Application) No. 2011-515194"}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic recording/reproduction apparatus, and more particularly, relates to a magnetic recording/reproduction apparatus provided with a mechanism for moving a sub-chassis relative to a main chassis.\n2. Description of the Related Art\nAn example of conventional magnetic recording/reproduction apparatus is disclosed in Japanese Patent Gazette No. 2627465. Referring to FIGS. 36 to 38, a conventional magnetic recording/reproduction apparatus 300 will be described.\nFIG. 36 is a plan view of a slide chassis 301. A spring hook 341 having a groove 356 is mounted on the slide chassis 301.\nFIG. 37 is a plan view of a main chassis 351. A slide chassis driving lever 359 is swingably mounted on the main chassis 351 via a swing shaft 361. A cam pin 363 extends from the slide chassis driving lever 359 and engages with a slide chassis driving groove 357 of the main cam 353.\nReferring to FIG. 38, the spring hook 341 has two holes 341A and 341B for fastening the spring hook 341 to the slide chassis 301. The holes 341A and 341B are elongate in the direction of sliding of the slide chassis 301 (direction indicated by arrow B3). Using these elongate holes 341A and 341B, the position of the spring hook 341 can be adjusted in the direction B3 so that the position of the slide chassis 301 and the position of a slide chassis driving pin 365 of the slide chassis driving lever 359 correspond to each other in the loading state. After the adjustment, the spring hook 341 is fastened to the slide chassis 301 with screws or the like. Such a construction is normal for a conventional magnetic recording/reproduction apparatus having a slide chassis.\nThe driving mechanism for the slide chassis 301 as described above has at least the following problems.\nIn the above mechanism, the position of the slide chassis 301 at the completion of the loading operation is adjusted by moving the slide chassis 301 in the direction B3 as described above. The spring hook 341 is fastened to the slide chassis 301 via the holes elongate in the direction B3. Accordingly, if an abnormal load is applied to the spring hook 341 in a direction blocking the loading/unloading operations during the loading/unloading operations, the spring hook 341 may move in the direction B3 against the fastening force, resulting in a change in the position of the slide chassis 301 with respect to the main chassis 351. Such an application of an abnormal load is likely to occur since the slide chassis is the portion of the magnetic recording/reproduction apparatus which is exposed to direct contact with a user.\nMoreover, a reel base for driving a tape reel in a cassette is located at a position in the direction B3 from the spring hook 341. Accordingly, an attempt to secure a moving space for the spring hook 341 will block realization of a small-sized mechanism.\nThe object of the present invention is to provide a magnetic recording/reproduction apparatus provided with a small-sized mechanism capable of achieving high reliability."}
{"text": "In Canada, cigarettes are packaged and sold in packets of 20 and 25 cigarettes per package, generally in two distinctly-different types of package, with each cigarette usually having a length of 85 mm or 100 mm. For the 20 cigarette-size package, a hard-box hinge-lid package encloses the cigarettes and a cover hinged to a lower cigarette-retaining portion is used for opening the package to allow access to the cigarettes and for reclosing the package. The cigarettes are arranged in a single bundle or group in three parallel rows, the outer rows containing seven cigarettes and the middle row containing six cigarettes, with each of the cigarettes in the middle row engaging two of the cigarettes in each of the outer rows.\nFor the 25 cigarette-size package, however, the cigarettes are supported in an inner tray which is slidably-mounted in an outer sleeve so that access to the cigarettes is obtained by sliding of the inner tray part way out of the sleeve. The cigarettes are separated into two distinct bundles or groups, one group containing 12 cigarettes arranged in two parallel rows and the outer group containing 13 cigarettes arranged in two blocks of two parallel rows of three cigarettes each, separated by the odd cigarette, the respective parallel rows of the groups being in straight line alignment.\nIn some instances, 20 cigarette-size packages of the same type as the 25 cigarette-size package mentioned above are used, and in this instance, the two bundles of cigarettes have 10 cigarettes each arranged in two parallel rows.\nOne of the drawbacks to the current 25-cigarette packages is their bulkiness due to the two row arrangement of the cigarettes in the package, leading to the necessity of shirt-pocket storage and transportation with the package on its side. This storage arrangement, however, is possible only with 85 mm or shorter cigarettes since shirt pockets do not have a width dimention to accommodate longer cigarettes. Further when the package contains 85 mm cigarettes, and is stored in this way, their is usually insufficient residual room to allow storage of writing implements or matches as well in the shirt pocket.\nFurthermore, storage and transportation of cigarettes on their side in this way leads to tobacco which has fallen out of the cigarettes being distributed over the cigarettes, leading commonly to tobacco particles on the outer surface of the cigarette filters, so that tobacco particles may enter the smoker's mouth unless the particles on the filter surface are carefully removed before the cigarette is placed in the smoker's mouth. Removal of these tobacco particles is a tedious chore for the smoker and entry of tobacco particles into the mouth is considered undesirable by many smokers.\nDespite these drawbacks, a considerable number of smokers prefer to purchase cigarettes in units of 25 rather than 20, even though the 20-cigarette package may be stored and transported upright in a shirt pocket, so that any tobacco particles falling out of the cigarettes remain in the bottom of the package, and hence the filter-fouling problem does not arise.\nAttempts to package 25 cigarettes in a hinge lid pack to take advantage of its unitary construction, compact form and lack of filter fouling by tobacco particles have not been successful since the increased width of package required to accommodate three rows of cigarettes in a single-bundle 25-cigarette array leads to the falling out of sight of cigarettes into the lower cigarette-retaining portion when only a few remain, giving the impression of an empty package and presenting difficulties in access to the cigarettes for removal from the package. This problem is more acute with the more-common 85 mm-length cigarettes as compared with the less-common 100 mm-length cigarettes. Where the tray-and-sleeve 25-cigarette package mentioned above is used, the accessibility problem does not arise, since simple sliding of the tray relative to the sleeve allows any remaining cigarette to be readily detected and removed. This package, however, has dimensional and other defects, as mentioned above.\nEven in the case of the 20-cigarette package, when one or a few cigarettes remain, typically of 85 mm or less length, there is a falling out of sight of the cigarettes. Although access to these cigarettes is a minor problem as compared with the more severe problem of a 25-cigarette package, it nevertheless exists but has been tolerated by the art.\nA further difficulty in packaging cigarettes in a three-row array in a hinge-lid pack arises from the need to have a cigarette located at each corner of the bundle for ease of wrapping of the bundle in foil paper to provide a cubic shape to the bundle. Thus, the outer rows of cigarettes must contain one more cigarette than the centre row and each cigarette in the centre row now must engage two cigarettes in each of the outer rows. These requirements allow only certain numbers of cigarettes to be provided in a single bundle, the number increasing by three for each increased size of bundle. The minimum number is five and the possible numbers of cigarettes include 20, 23 and 26, but not 25.\nThus, heretofore, there has never been provided a cigarette package of the hinge-lid type containing total numbers of cigarettes of 20 or more arranged in three rows and which allows all the cigarettes in the package to be visible and accessible irrespective of the number of cigarettes remaining in the package."}
{"text": "The present invention relates to an aqueous, preferably clear, composition, article of manufacture, and method for use as a freshening composition. Preferably, the composition is sprayed onto fabrics, particularly clothes, to restore their freshness by reducing malodor impression, without washing or dry cleaning. Fabrics treated with the composition of the present invention also release extra fragrance upon rewetting, such as when the wearer perspires. The freshening composition of the present invention is designed to extend the wear of fabrics between washing or dry cleaning. Fabrics treated with the freshening composition of the present invention will stay fresher longer, and receive extra freshening effect via perfume release when it is most needed, that is upon fabric rewetting.\nA wide variety of deodorizing compositions are known in the art, the most common of which only contain a perfume to mask the malodor. Odor masking is the intentional concealment of one odor by the addition of another. The preference to the masking perfume is varied greatly, depending on the application, e.g., underarm odor masking, fabric odor masking, bathroom odor masking, etc. Appropriate perfume ingredients need to be selected to connote freshness.\nOdor modification, in which the odor is changed, e.g., by chemical modification, has also been used. Current malodor modification methods known in the art which do not simply mask odors are oxidative degradation, which uses oxidizing agents such as oxygen bleaches, chlorine, chlorinated materials such as sodium hypochlorite, chlorine dioxide, etc., and potassium permanganate to reduce malodor, and reductive degradation which uses reducing agents such as sodium bisulfite to reduce malodor. Both of these methods are unacceptable for use on fabrics because they can damage colored fabrics, specifically, they can bleach and discolor colored fabrics.\nOther methods of odor control contain actives that are targeted to react with malodors having specific chemical functional groups. Examples of such actives are: biguanide polymers, which complex with organic compounds containing organically bound N and/or S atoms and fatty alcohol esters of methyl methacrylic acid which react with thiols, amines, and aldehydes. A more detailed description of these methods can be found in U.S. Pat. Nos. 2,544,093, 3,074,891, and U.K. Pat. App. No. 941,105, all of said patents and applications are incorporated herein by reference. Fatty alcohol esters of methyl methacrylic acid are not preferred in the composition of this invention because they are not water soluble.\nOther types of deodorizing compositions known in the art contain antibacterial and antifungal agents which regulate the malodor-producing microorganisms found on the surface to which the deodorizing composition is directed. Many skin deodorant products use this technology. These compositions are not effective on malodors that do not come from bacterial sources, such as tobacco or food odors.\nFabric malodor is most commonly caused by environmental odors such as tobacco odor, cooking and/or food odors, or body odor. The unpleasant odors are mainly organic molecules which have different structures and functional groups. One type of malodor that is very noticeable, and is commonly found on worn fabrics is low molecular weight, straight-chain, branched, and unsaturated C.sub.6 -C.sub.11 fatty acids that cause axillary odor. See \"Analysis of Characteristic Odor from Human Male Axillae\", X. Zeng, et al., J. Chem. Ecol., pp. 1469-1492, 1991, incorporated herein by reference. See also, U.S. Pat. No. 4,664,909, Marschner et al., issued May 12, 1987, BE 830,098, published Oct. 1, 1975, and CA 1,088,428, published Oct. 28, 1980, DE 2,803,176, published Aug. 3, 1978, all of said patents and applications incorporated herein by reference.\nCyclodextrin molecules are known for their ability to form complexes with perfume ingredients and have typically been taught as a perfume carrier. The prior art teaches the use of drier-added fabric softener sheets containing high levels of cyclodextrin/perfume complexes wherein the fabrics treated with this solid cyclodextrin complex release perfume when the fabrics are rewetted. The art also teaches that cyclodextrin/perfume complexes used in aqueous rinse-added fabric softener compositions must be protected, e.g., with a hydrophobic wax coating so the cyclodextrin/perfume complexes will not decompose due to the presence of water. See U.S. Pat. No. 5,102,564 Gardlik et al., issued Apr. 7, 1992; U.S. Pat. No. 5,234,610, Gardlik et al., issued Aug. 10, 1993; U.S. Pat. No. 5,234,611 Trinh, et al., issued Aug. 10, 1993, all of said patents incorporated herein by reference. It is therefore highly surprising and unexpected to find that fabrics treated with the aqueous compositions of the present invention, which contain low levels of cyclodextrin, also exhibit perfume release upon rewetting. This phenomenon creates a benefit in that fabrics treated with the composition of the present invention will thus remain fresh longer, via a perfume release, when said fabrics are rewetted, such as when the wearer perspires."}
{"text": "In the recent significant development of smaller-size and lighter-weight optical equipment, aspherical lenses have been used increasingly. The aspherical lens is advantageous in that aberration of light can readily be corrected and that the number of lenses can be decreased so as to allow reduction in size of the equipment.\nFor fabricating an aspherical lens or the like, a glass preform is softened by heating, which is then formed into a desired shape by precision-mold press molding. There are generally two ways of obtaining the preform: one is to cut a piece of glass out of a glass block or bar and process it into a preform, and the other is to drop a molten glass from a distal end of a nozzle so as to obtain a glass preform in the spherical form.\nIn order to obtain a molded product of glass by way of precision molding, it is necessary to press-mold the preform under the temperature condition near the deformation point (At). Therefore, when the preform has a higher deformation point (At), the mold coming into contact with the preform will be exposed to a higher temperature, causing the surface of the mold to suffer oxidization and corrosion. This gives rise to the need of maintenance of the mold, hindering mass production at a low cost. Accordingly, it is desired that the optical glass constituting the preform can be molded at a relatively low temperature, or, that it has a low glass transition point (Tg) and/or a low deformation point (At).\nAs to the glass used for a molded lens, a glass having various optical characteristics suitable for its specific use is demanded. In particular, there is an increasing demand for a glass having a high refractive index, low dispersion, and a low deformation point.\nThe conventional glasses satisfying the above-described optical characteristics include a barium flint glass. This not only contains PbO (lead oxide) hazardous to humans, but also poses other unfavorable problems. For example, metallic lead would be deposited on a surface of the product upon precision press molding, and a glass surface would be likely to become rough due to fusion with the mold.\nIn a digital camera, it is necessary to reduce the lens surface reflection as much as possible, and anti-reflection coating is used for that purpose. In order to restrict the reflectance as well as incident angle dependence and also to broaden the wavelength band, however, a considerable number of layers of coating films are required, resulting in complicated and expensive process steps.\nAs a way of achieving low reflectance without the coating films, it is known to form, on the surface of a lens or the like, a fine structure that is smaller in size than the wavelength of light. This may be done, for example, by nanoimprinting using a resin. A material having a low softening temperature such as a resin is relatively easy to form and shape using a microfabricated mold. However, temperature dependence of refractive index of the resin is about −1×10−4 (K), which is greater than that of the glass by two orders of magnitude. This means that for a part intended for higher image quality, the change in refractive index will affect the image quality more severely. In view of the foregoing, a glass for use in transferring a fine structure has been studied in order to achieve higher functionality of an optical part. A fine structure-transferred glass is an ultra-precision-mold press-molded product having a glass surface onto which the mold's concavo-convex pattern on the order of μm to nm has been transferred. For example, a conventional optical part may be replaced with a lens provided with such a fine structure so as to advantageously achieve a compact device with higher functionality at a reduced cost. A mold made up of Ni and P, which has conventionally been used primarily for molding a resin lens, may be used as well. As the characteristics required for the glass, it is crucially important that the deformation point is 500° C. or lower in order to restrict deterioration of the mold. Furthermore, in order to eliminate the need of a coating film on a lens, climate resistance of the glass itself is required as well.\nAs a glass free of PbO and having the above-described optical characteristics, a P2O5—R12O—R2O-(rare earth oxide or the like) type glass (where R1: alkali metal oxide, and R2: divalent metal oxide) has been disclosed. This optical glass has a refractive index (nd) of 1.63 to 1.67, an Abbe number (υd) of 47 to 59, and a deformation point (At) of 500° C. or lower (Patent Document 1).\nThere is also disclosed a P2O5—R12O—BaO—ZnO-(high-valent oxide) type optical glass. This optical glass has a refractive index (nd) of 1.52 to 1.7, and an Abbe number (υd) of 42 to 70 (Patent Document 2).\nAlso disclosed are a P2O5—R12O—R2O—Nb2O5 type glass, a P2O5—R12O—Nb2O5—WO3 type glass, and a P2O5—R12O—Bi2O3 type glass. These optical glasses each have a refractive index (nd) of 1.57 or greater, and a deformation point (At) of 570° C. or lower (Patent Documents 3 to 7).\nThere are also disclosed a P2O5—B2O3—R12O—R2O—Gd2O3 type glass and a P2O5—B2O3—R12O—BaO—ZnO type glass. These optical glasses each have a refractive index (nd) of L54 or greater, and an Abbe number (υd) of 57 or greater (Patent Documents 8 and 9).\nFurther disclosed is a P2O5—R12O—R2O—ZnO—Al2O3 type glass. This optical glass has a refractive index (nd) of 1.55 to 1.65, an Abbe number (υd) of 55 to 65, and a deformation point of 500° C. or lower (Patent Documents 10 to 12).\nPatent Document 1: Japanese Patent Application Laid-Open No. 11-139845\nPatent Document 2: Japanese Patent Application Laid-Open No. 2004-217513\nPatent Document 3: Japanese Patent Application Laid-Open No. 2002-293572\nPatent Document 4: Japanese Patent Application Laid-Open No. 2005-247659\nPatent Document 5: Japanese Patent Application Laid-Open No. 2003-335549\nPatent Document 6: Japanese Patent Application Laid-Open No. 2004-2153\nPatent Document 7: Japanese Patent Application Laid-Open No. 2003-238197\nPatent Document 8: Japanese Patent Application Laid-Open No. 2006-52119\nPatent Document 9: WO 2003/072518\nPatent Document 10: Japanese Patent Application Laid-Open No. 2004-168593\nPatent Document 11: Japanese Patent Application Laid-Open No. 2004-262703\nPatent Document 12: Japanese Patent Application Laid-Open No. 2005-53749"}
{"text": "1. Field of the Invention\nThe present invention relates generally to automatic electronic flash, and particularly concerns the automatic electronic flash capable of automatic adjustment of amount of flashed light when photographic scenery is of rear light.\n2. Description of the Prior Art\nAutomatic electronic flash, which has a function of automatically controlling amount of flashed light in response to the light reflected from a scene, is very useful to make good picture for an amateur photographer.\nIn an automatic electronic flash photographing for a rear light scene for instance in outdoor scene with portrait in a bright sky, namely in a fill-in flash photographing, there are the following problems.\nWhen the automatic electronic photographing for the rear light scene is executed on trial, most probable operation sequence is: firstly, brightness of scene is measured by exposure measuring system of a camera for the photographic scene angle for a shutter speed appropriate for strobe synchronization, thereby to obtain suitable F-stop value for the scene, and secondly, F-stop number of the automatic electronic flash and the stop of the camera are adjusted to the above-mentioned suitable stop value.\nWhen a rear light scene is considered, in most cases the back scene of the object is a vacant space, or even when not the vacant space, something to reflect light is at a very far distance. Accordingly, unlike the ordinary indoor flash photographing with automatic electronic flash, amount of light incident to light sensor of the automatic electronic flash during the photographic flash is very small. In other word, unlike the indoor photographing where lights are reflected from various matters surrounding the photographic object, in case of a rear light photographing, for instance in outdoor scene, there is substantially no reflection of light of the flash, since there are no surrounding matters in not far distance other than the photographic object. As a result, in the rear light photographing, the amount of light incident to light sensor of the automatic electronic flash becomes considerably smaller than that in the case of indoor flash phographing with the same automatic electronic flash.\nAccordingly, when the automatic photographing of outdoor is carried out in the above-mentioned operation sequence, the photographing naturally results in an overexposuring, thereby making the resultant photograph whitish.\nAccording to the experimental study of the inventors, it has been found that, for fill-in flash photographing for such scenes with mountain or sea as backgound, the photograph becomes overexposured by a degree of 1-2 F-stops as converted into the stop value, irrespective of the stop value used within the usually used stop value range of F1.4-F22.\nAccordingly, it has been believed that the fill-in flash photographing is not satisfactorily made with automatic electronic flashing, and instead a manual flash apparatus has been used as follows: firstly shutter speed of the camera is set usually to 1/60 sec. and the brightness of the photographic scene is measured by appropriately setting film sensitivity in the camera or in an exposure meter, thereby to measure suitable stop value. And secondly, distance between the camera and the photographic object is measured and a product of the distance and the above-mentioned stop value is calculated and a guide number dial of the manual flash apparatus is adjusted to the value of the product.\nThe above-mentioned conventional manual setting of the flash apparatus, however, is troublesome for amature photographer of small experience, and the above-mentioned sequence of manual setting is not well understood, or setting of the guide number by the product is difficult for him or been forgot by him. Accordingly there are many liability of failure of the flash photography in backlight. And therefore, the rear light photography is considered difficult to produce good picture.\nFurthermore in some case, even when the user can understand and operates the above-mentioned sequence of the manual flash photographing, it is necessary that the calculated guide number is always available on the flash apparatus, and in such case there is a necessity to change distance between the camera and the object or in some time the photographing of the object becomes impossible.\nOn the other hand, in order to solve the above-mentioned inconvenience, there has been many proposals for automatic electronic flashes easily usable for amateur photographer, but these modern devices are limited in a parmanent combination with camera. In view of the above-mentioned problem, there is a strong demand for a satisfactory automatic electronic flash capable of photographing even for rear light scenery without troublesome manual operation and calculation."}
{"text": "The detection and characterization of specific nucleic acid sequences and sequence changes have been utilized to detect the presence of viral or bacterial nucleic acid sequences indicative of an infection, the presence of variants or alleles of mammalian genes associated with disease and cancers, and the identification of the source of nucleic acids found in forensic samples, as well as in paternity determinations.\nVarious methods are known in the art which may be used to detect and characterize specific nucleic acid sequences and sequence changes. Nonetheless, as nucleic acid sequence data of the human genome, as well as the genomes of pathogenic organisms accumulates, the demand for fast, reliable, cost-effective and user-friendly tests for specific sequences continues to grow. Importantly, these tests must be able to create a detectable signal from a very low copy number of the sequence of interest. The following discussion examines three levels of nucleic acid detection currently in use: I. Signal Amplification Technology for detection of rare sequences; II. Direct Detection Technology for detection of higher copy number sequences; and III. Detection of Unknown Sequence Changes for rapid screening of sequence changes anywhere within a defined DNA fragment.\nI. Signal Amplification Technology Methods For Amplification\nThe \"Polymerase Chain Reaction\" (PCR) comprises the first generation of methods for nucleic acid amplification. However, several other methods have been developed that employ the same basis of specificity, but create signal by different amplification mechanisms. These methods include the \"Ligase Chain Reaction\" (LCR), \"Self-Sustained Synthetic Reaction\" (3SR/NASBA), and \"Q.beta.-Replicase\" (Q.beta.).\nPolymerase Chain Reaction (PCR)\nThe polymerase chain reaction (PCR), as described in U.S. Pat. Nos. 4,683,195 and 4,683,202 to Mullis and Mullis et al., describe a method for increasing the concentration of a segment of target sequence in a mixture of genomic DNA without cloning or purification. This technology provides one approach to the problems of low target sequence concentration. PCR can be used to directly increase the concentration of the target to an easily detectable level. This process for amplifying the target sequence involves introducing a molar excess of two oligonucleotide primers which are complementary to their respective strands of the double-stranded target sequence to the DNA mixture containing the desired target sequence. The mixture is denatured and then allowed to hybridize. Following hybridization, the primers are extended with polymerase so as to form complementary strands. The steps of denaturation, hybridization, and polymerase extension can be repeated as often as needed, in order to obtain relatively high concentrations of a segment of the desired target sequence.\nThe length of the segment of the desired target sequence is determined by the relative positions of the primers with respect to each other, and, therefore, this length is a controllable parameter. Because the desired segments of the target sequence become the dominant sequences (in terms of concentration) in the mixture, they are said to be \"PCR-amplified.\"\nLigase Chain Reaction (LCR or LAR)\nThe ligase chain reaction (LCR; sometimes referred to as \"Ligase Amplification Reaction\" (LAR) described by Barany, Proc. Natl. Acad. Sci., 88:189 (1991); Barany, PCR Methods and Applic., 1:5 (1991); and Wu and Wallace, Genomics 4:560 (1989) has developed into a well-recognized alternative method for amplifying nucleic acids. In LCR, four oligonucleotides, two adjacent oligonucleotides which uniquely hybridize to one strand of target DNA, and a complementary set of adjacent oligonucleotides, which hybridize to the opposite strand are mixed and DNA ligase is added to the mixture. Provided that there is complete complementarity at the junction, ligase will covalently link each set of hybridized molecules. Importantly, in LCR, two probes are ligated together only when they base-pair with sequences in the target sample, without gaps or mismatches. Repeated cycles of denaturation, hybridization and ligation amplify a short segment of DNA. LCR has also been used in combination with PCR to achieve enhanced detection of single-base changes. Segev, PCT Public. No. W09001069 A1 (1990). However, because the four oligonucleotides used in this assay can pair to form two short ligatable fragments, there is the potential for the generation of target-independent background signal. The use of LCR for mutant screening is limited to the examination of specific nucleic acid positions.\nSelf-Sustained Synthetic Reaction (3SR/NASBA)\nThe self-sustained sequence replication reaction (3SR) (Guatelli et al., Proc. Natl. Acad. Sci., 87:1874-1878 1990!, with an erratum at Proc. Natl. Acad. Sci., 87:7797 1990!) is a transcription-based in vitro amplification system (Kwok et al., Proc. Natl. Acad. Sci., 86:1173-1177 1989!) that can exponentially amplify RNA sequences at a uniform temperature. The amplified RNA can then be utilized for mutation detection (Fahy et al., PCR Meth. Appl., 1:25-33 1991!). In this method, an oligonucleotide primer is used to add a phage RNA polymerase promoter to the 5' end of the sequence of interest. In a cocktail of enzymes and substrates that includes a second primer, reverse transcriptase, RNase H, RNA polymerase and ribo-and deoxyribonucleoside triphosphates, the target sequence undergoes repeated rounds of transcription, cDNA synthesis and second-strand synthesis to amplify the area of interest. The use of 3SR to detect mutations is kinetically limited to screening small segments of DNA (e.g., 200-300 base pairs).\nQ-Beta (Q.beta.) Replicase\nIn this method, a probe which recognizes the sequence of interest is attached to the replicatable RNA template for Q.beta. replicase. A previously identified major problem with false positives resulting from the replication of unhybridized probes has been addressed through use of a sequence-specific ligation step. However, available thermostable DNA ligases are not effective on this RNA substrate, so the ligation must be performed by T4 DNA ligase at low temperatures (37.degree. C.). This prevents the use of high temperature as a means of achieving specificity as in the LCR, the ligation event can be used to detect a mutation at the junction site, but not elsewhere.\nTable 1 below, lists some of the features desirable for systems useful in sensitive nucleic acid diagnostics, and summarizes the abilities of each of the major amplification methods (See also, Landgren, Trends in Genetics 9:199 1993!).\nA successful diagnostic method must be very specific. A straight-forward method of controlling the specificity of nucleic acid hybridization is by controlling the temperature of the reaction. While the 3SR/NASBA, and Q.beta. systems are all able to generate a large quantity of signal, one or more of the enzymes involved in each cannot be used at high temperature (i.e., &gt;55.degree. C.). Therefore the reaction temperatures cannot be raised to prevent non-specific hybridization of the probes. If probes are shortened in order to make them melt more easily at low temperatures, the likelihood of having more than one perfect match in a complex genome increases. For these reasons, PCR and LCR currently dominate the research field in detection technologies.\nTABLE 1 ______________________________________ METHOD: PCR & 3SR FEATURE PCR LCR LCR NASBA Q.beta. ______________________________________ Amplifies Target + + + + Recognition of Independent + + + + + Sequences Required Performed at High Temp. + + Operates at Fixed Temp. + + Exponential Amplification + + + + + Generic Signal Generation + Easily Automatable ______________________________________\nThe basis of the amplification procedure in the PCR and LCR is the fact that the products of one cycle become usable templates in all subsequent cycles, consequently doubling the population with each cycle. The final yield of any such doubling system can be expressed as: (1+X).sup.n =y, where \"X\" is the mean efficiency (percent copied in each cycle), \"n\" is the number of cycles, and \"y\" is the overall efficiency, or yield of the reaction (Mullis, PCR Methods Applic., 1:1 1991!). If every copy of a target DNA is utilized as a template in every cycle of a polymerase chain reaction, then the mean efficiency is 100%. If 20 cycles of PCR are performed, then the yield will be 2.sup.20, or 1,048,576 copies of the starting material. If the reaction conditions reduce the mean efficiency to 85%, then the yield in those 20 cycles will be only 1.85.sup.20, or 220,513 copies of the starting material. In other words, a PCR running at 85% efficiency will yield only 21% as much final product, compared to a reaction running at 100% efficiency. A reaction that is reduced to 50% mean efficiency will yield less than 1% of the possible product.\nIn practice, routine polymerase chain reactions rarely achieve the theoretical maximum yield, and PCRs are usually run for more than 20 cycles to compensate for the lower yield. At 50% mean efficiency, it would take 34 cycles to achieve the million-fold amplification theoretically possible in 20, and at lower efficiencies, the number of cycles required becomes prohibitive. In addition, any background products that amplify with a better mean efficiency than the intended target will become the dominant products.\nAlso, many variables can influence the mean efficiency of PCR, including target DNA length and secondary structure, primer length and design, primer and dNTP concentrations, and buffer composition, to name but a few. Contamination of the reaction with exogenous DNA (e.g., DNA spilled onto lab surfaces) or cross-contamination is also a major consideration. Reaction conditions must be carefully optimized for each different primer pair and target sequence, and the process can take days, even for an experienced investigator. The laboriousness of this process, including numerous technical considerations and other factors, presents a significant drawback to using PCR in the clinical setting. Indeed, PCR has yet to penetrate the clinical market in a significant way. The same concerns arise with LCR, as LCR must also be optimized to use different oligonucleotide sequences for each target sequence. In addition, both methods require expensive equipment, capable of precise temperature cycling.\nMany applications of nucleic acid detection technologies, such as in studies of allelic variation, involve not only detection of a specific sequence in a complex background, but also the discrimination between sequences with few, or single, nucleotide differences. One method for the detection of allele-specific variants by PCR is based upon the fact that it is difficult for Taq polymerase to synthesize a DNA strand when there is a mismatch between the template strand and the 3' end of the primer. An allele-specific variant may be detected by the use of a primer that is perfectly matched with only one of the possible alleles; the mismatch to the other allele acts to prevent the extension of the primer, thereby preventing the amplification of that sequence. This method has a substantial limitation in that the base composition of the mismatch influences the ability to prevent extension across the mismatch, and certain mismatches do not prevent extension or have only a minimal effect (Kwok et al., Nucl. Acids Res., 18:999 1990!).)\nA similar 3'-mismatch strategy is used with greater effect to prevent ligation in the LCR (Barany, PCR Meth. Applic., 1:5 1991!). Any mismatch effectively blocks the action of the thermostable ligase, but LCR still has the drawback of target-independent background ligation products initiating the amplification. Moreover, the combination of PCR with subsequent LCR to identify the nucleotides at individual positions is also a clearly cumbersome proposition for the clinical laboratory.\nII. Direct Detection Technology\nWhen a sufficient amount of a nucleic acid to be detected is available, there are advantages to detecting that sequence directly, instead of making more copies of that target, (e.g., as in PCR and LCR). Most notably, a method that does not amplify the signal exponentially is more amenable to quantitative analysis. Even if the signal is enhanced by attaching multiple dyes to a single oligonucleotide, the correlation between the final signal intensity and amount of target is direct. Such a system has an additional advantage that the products of the reaction will not themselves promote further reaction, so contamination of lab surfaces by the products is not as much of a concern. Traditional methods of direct detection including Northern and Southern blotting and RNase protection assays usually require the use of radioactivity and are not amenable to automation. Recently devised techniques have sought to eliminate the use of radioactivity and/or improve the sensitivity in automatable formats. Two examples are the \"Cycling Probe Reaction\" (CPR), and \"Branched DNA\" (bDNA)\nThe cycling probe reaction (CPR) (Duck et al., BioTech., 9:142 1990!), uses a long chimeric oligonucleotide in which a central portion is made of RNA while the two termini are made of DNA. Hybridization of the probe to a target DNA and exposure to a thermostable RNase H causes the RNA portion to be digested. This destabilizes the remaining DNA portions of the duplex, releasing the remainder of the probe from the target DNA and allowing another probe molecule to repeat the process. The signal, in the form of cleaved probe molecules, accumulates at a linear rate. While the repeating process increases the signal, the RNA portion of the oligonucleotide is vulnerable to RNases that may carried through sample preparation.\nBranched DNA (bDNA), described by Urdea et al., Gene 61:253-264 (1987), involves oligonucleotides with branched structures that allow each individual oligonucleotide to carry 35 to 40 labels (e.g., alkaline phosphatase enzymes). While this enhances the signal from a hybridization event, signal from non-specific binding is similarly increased.\nIII. Detection Of Unknown Sequence Changes\nThe demand for tests which allow the detection of specific nucleic acid sequences and sequence changes is growing rapidly in clinical diagnostics. As nucleic acid sequence data for genes from humans and pathogenic organisms accumulates, the demand for fast, cost-effective, and easy-to-use tests for as yet unknown mutations within specific sequences is rapidly increasing.\nA handful of methods have been devised to scan nucleic acid segments for mutations. One option is to determine the entire gene sequence of each test sample (e.g., a bacterial isolate). For sequences under approximately 600 nucleotides, this may be accomplished using amplified material (e.g., PCR reaction products). This avoids the time and expense associated with cloning the segment of interest. However, specialized equipment and highly trained personnel are required, and the method is too labor-intense and expensive to be practical and effective in the clinical setting.\nIn view of the difficulties associated with sequencing, a given segment of nucleic acid may be characterized on several other levels. At the lowest resolution, the size of the molecule can be determined by electrophoresis by comparison to a known standard run on the same gel. A more detailed picture of the molecule may be achieved by cleavage with combinations of restriction enzymes prior to electrophoresis, to allow construction of an ordered map. The presence of specific sequences within the fragment can be detected by hybridization of a labeled probe, or the precise nucleotide sequence can be determined by partial chemical degradation or by primer extension in the presence of chain-terminating nucleotide analogs.\nFor detection of single-base differences between like sequences, the requirements of the analysis are often at the highest level of resolution. For cases in which the position of the nucleotide in question is known in advance, several methods have been developed for examining single base changes without direct sequencing. For example, if a mutation of interest happens to fall within a restriction recognition sequence, a change in the pattern of digestion can be used as a diagnostic tool (e.g., restriction fragment length polymorphism RFLP! analysis).\nSingle point mutations have been also detected by the creation or destruction of RFLPs. Mutations are detected and localized by the presence and size of the RNA fragments generated by cleavage at the mismatches. Single nucleotide mismatches in DNA heteroduplexes are also recognized and cleaved by some chemicals, providing an alternative strategy to detect single base substitutions, generically named the \"Mismatch Chemical Cleavage\" (MCC) (Gogos et al., Nucl. Acids Res., 18:6807-6817 1990!). However, this method requires the use of osmium tetroxide and piperidine, two highly noxious chemicals which are not suited for use in a clinical laboratory.\nRFLP analysis suffers from low sensitivity and requires a large amount of sample. When RFLP analysis is used for the detection of point mutations, it is, by its nature, limited to the detection of only those single base changes which fall within a restriction sequence of a known restriction endonuclease. Moreover, the majority of the available enzymes have 4 to 6 base-pair recognition sequences, and cleave too frequently for many large-scale DNA manipulations (Eckstein and Lilley (eds.), Nucleic Acids and Molecular Biology, vol. 2, Springer-Verlag, Heidelberg 1988!). Thus, it is applicable only in a small fraction of cases, as most mutations do not fall within such sites.\nA handful of rare-cutting restriction enzymes with 8 base-pair specificities have been isolated and these are widely used in genetic mapping, but these enzymes are few in number, are limited to the recognition of G+C-rich sequences, and cleave at sites that tend to be highly clustered (Barlow and Lehrach, Trends Genet., 3:167 1987!). Recently, endonucleases encoded by group I introns have been discovered that might have greater than 12 base-pair specificity (Perlman and Butow, Science 246:1106 1989!), but again, these are few in number.\nIf the change is not in a recognition sequence, then allele-specific oligonucleotides (ASOs), can be designed to hybridize in proximity to the unknown nucleotide, such that a primer extension or ligation event can be used as the indicator of a match or a mis-match. Hybridization with radioactively labeled allelic specific oligonucleotides (ASO) also has been applied to the detection of specific point mutations (Conner et al., Proc. Natl. Acad. Sci., 80:278-282 1983!). The method is based on the differences in the melting temperature of short DNA fragments differing by a single nucleotide. Stringent hybridization and washing conditions can differentiate between mutant and wild-type alleles. The ASO approach applied to PCR products also has been extensively utilized by various researchers to detect and characterize point mutations in ras genes (Vogelstein et al., N. Eng. J. Med., 319:525-532 1988!; and Farr et al., Proc. Natl. Acad. Sci., 85:1629-1633 1988!), and gsp/gip oncogenes (Lyons et al., Science 249:655-659 1990!). Because of the presence of various nucleotide changes in multiple positions, the ASO method requires the use of many oligonucleotides to cover all possible oncogenic mutations.\nWith either of the techniques described above (i.e., RFLP and ASO), the precise location of the suspected mutation must be known in advance of the test. That is to say, they are inapplicable when one needs to detect the presence of a mutation of an unknown character and position within a gene or sequence of interest.\nTwo other methods rely on detecting changes in electrophoretic mobility in response to minor sequence changes. One of these methods, termed \"Denaturing Gradient Gel Electrophoresis\" (DGGE) is based on the observation that slightly different sequences will display different patterns of local melting when electrophoretically resolved on a gradient gel. In this manner, variants can be distinguished, as differences in melting properties of homoduplexes versus heteroduplexes differing in a single nucleotide can detect the presence of mutations in the target sequences because of the corresponding changes in their electrophoretic mobilities. The fragments to be analyzed, usually PCR products, are \"clamped\" at one end by a long stretch of G-C base pairs (30-80) to allow complete denaturation of the sequence of interest without complete dissociation of the strands. The attachment of a GC \"clamp\" to the DNA fragments increases the fraction of mutations that can be recognized by DGGE (Abrams et al., Genomics 7:463-475 1990!). Attaching a GC clamp to one primer is critical to ensure that the amplified sequence has a low dissociation temperature (Sheffield et al., Proc. Natl. Acad. Sci., 86:232-236 1989!; and Lerman and Silverstein, Meth. Enzymol., 155:482-501 1987!). Modifications of the technique have been developed, using temperature gradients (Wartell et al., Nucl. Acids Res., 18:2699-2701 1990!), and the method can be also applied to RNA:RNA duplexes (Smith et al., Genomics 3:217-223 1988!).\nLimitations on the utility of DGGE include the requirement that the denaturing conditions must be optimized for each type of DNA to be tested. Furthermore, the method requires specialized equipment to prepare the gels and maintain the needed high temperatures during electrophoresis. The expense associated with the synthesis of the clamping tail on one oligonucleotide for each sequence to be tested is also a major consideration.\nAnother common method, called \"Single-Strand Conformation Polymorphism\" (SSCP) was developed by Hayashi, Sekya and colleagues (reviewed by Hayashi, PCR Meth. Appl., 1:34-38, 1991!) and is based on the observation that single strands of nucleic acid can take on characteristic conformations in non-denaturing conditions, and these conformations influence electrophoretic mobility. The complementary strands assume sufficiently different structures that one strand may be resolved from the other. Changes in sequences within the fragment will also change the conformation, consequently altering the mobility and allowing this to be used as an assay for sequence variations (Orita, et al., Genomics 5:874-879, 1989!).\nThe SSCP process involves denaturing a DNA segment (e.g., a PCR product) that is labelled on both strands, followed by slow electrophoretic separation on a non-denaturing polyacrylamide gel, so that intra-molecular interactions can form and not be disturbed during the run. This technique is extremely sensitive to variations in gel composition and temperature. A serious limitation of this method is the relative difficulty encountered in comparing data generated in different laboratories, under apparently similar conditions.\nIn addition to the above limitations, all of these methods are limited as to the size of the nucleic acid fragment that can be analyzed. For the direct sequencing approach, sequences of greater than 600 base pairs require cloning, with the consequent delays and expense of either deletion sub-cloning or primer walking, in order to cover the entire fragment. SSCP and DGGE have even more severe size limitations. Because of reduced sensitivity to sequence changes, these methods are not considered suitable for larger fragments. Although SSCP is reportedly able to detect 90% of single-base substitutions within a 200 base-pair fragment, the detection drops to less than 50% for 400 base pair fragments. Similarly, the sensitivity of DGGE decreases as the length of the fragment reaches 500 base-pairs.\nClearly, there remains a need for a method that is less sensitive to size so that entire genes, rather than gene fragments, may be analyzed. Such a tool must also be robust, so that data from different labs, generated by researchers of diverse backgrounds and skills will be comparable. Ideally, such a method would be compatible with \"multiplexing,\" (i.e., the simultaneous analysis of several molecules or genes in a single reaction or gel lane, usually resolved from each other by differential labelling or probing). Such an analytical procedure would facilitate the use of internal standards for subsequent analysis and data comparison, and increase the productivity of personnel and equipment. The ideal method would also be easily automatable."}
{"text": "The class of polymers of carbon monoxide and olefin(s) is well known in the art. Brubaker, U.S. Pat. No. 2,495,286 produced such polymers of relatively low carbon monoxide content in the presence of free radical initiators, e.g., peroxy compounds. G.B. 1,081,304 produced similar polymers of higher carbon monoxide content in the presence of alkylphosphine complexes of palladium compounds as catalyst. Nozaki extended the reaction to produce linear alternating polymers in the presence of arylphosphine complexes of palladium moieties and certain inert solvents. See, for example, U.S. Pat. No. 3,694,412. More recently, the production of the linear alternating polymers has become of greater interest because of the availability and the desirable properties of such polymers, now becoming known as polyketones or polyketone polymers. The more recent processes for the production of polyketone polymers are illustrated by a number of published European patent applications including 121,965, 181,014, 213,671 and 257,663. The processes generally involve the use of a catalyst composition formed from a compound of palladium, cobalt or nickel, the anion of a strong non-hydrohalogenic acid and a bidentate ligand of phosphorus, arsenic or antimony. The scope of the polymerization is extensive but, without wishing to be limited, many preferred catalyst compositions are produced from a compound of palladium, the anion of a non-hydrohalogenic acid having a pKa below 2 and a bidentate aromatic ligand of phosphorus.\nThe rate at which polymerization takes place will be determined in part by the ethylenically unsaturated hydrocarbons undergoing polymerization. In general, ethylene will polymerize at a rate considerably faster than ethylenically unsaturated hydrocarbons of three or more carbon atoms. Copolymers of ethylene and a second ethylenically unsaturated hydrocarbon of three or more carbon atoms are also formed at an acceptable rate in the presence of the above-described catalyst composition. When all the ethylenically unsaturated hydrocarbons are of three or more carbon atoms the rate of polymerization decreases.\nIt is known that the rate of production of linear alternating polymers containing at least a portion of ethylene can be improved or promoted by the presence of organic oxidizing agents including hydroquinones. In U.S. Pat. No. 4,810,774 and U.S. Pat. No. 4,824,935, it is disclosed that oxygen-containing generally acyclic aliphatic compounds such as ethers, ketones and esters will serve to promote catalytic activity in the polymerization or copolymerization of ethylene to form linear alternating polymers. It would be of advantage to provide a process for the production of linear alternating polymers of carbon monoxide and ethylenically unsaturated hydrocarbon of at least 3 carbon atoms wherein the catalytic activity of a palladium-containing polyketone polymerization catalyst has been promoted."}
{"text": "According to the completion of the genome project for both human and various target animals and plants and the development of bioinformatics, mRNA has been proved to act as a messenger transmitting genetic information of DNA to a protein and at the same time to regulate the gene expression.\nSince the beginning of year 2000, micro-RNA (miRNA) or its precursor pre-miRNA has been proved to regulate 10-20% of gene functions. In prokaryotes, some parts of mRNA are directly bound with a metabolite, suggesting that it has ribo-switch that regulates the functions of metabolite related protein. It has been also confirmed that the secondary structure of untranslating region of mRNA of higher animals regulates mRNA stability and translation efficiency.\nThe numbers of such RNA that has regulatory function are considerable. The structure of the RNA is composed of a series of hairpin structures in which stems and loops (basic motif) are arranged serially. It is also presumed that pharmacophore of natural miRNA or biologically significant mRNA might be the specific stem-loop (hairpin) structure, which is less than 30 nt, considering the size of binding region of ribo-switch to a compound.\nAlthough every mRNA has been proved to have secondary structure, the confirmed mRNA hairpin structures are very few, which are only exemplified by Rev Response Element (RRE) of HIV-1, trans-activation response element (TAR) of HIV-1, Thymidylate Synthase mRNA of various tumor cells and Ion Responsive Element (IRE) involved in homeostasis of iron ion and dementia, which is attributed to the lack of biological methods, the lack of information on RNA-binding protein and insufficient information on hairpin structure, etc. However, RNA targets having the hairpin structure are highly expected to be major biological targets and so great effort has to be made to find out ligands against such pouring RNA targets.\nPolyamines having several amine groups have been produced by imitating the conventional RNA pro-binding aminoglycoside compound, which have also been confirmed to be bound with RNA targets very well (Lawton et al., J. Am. Chem. Soc., 126: 12762-12763, 2004). Successively, morphology of a protein was observed according to the methylation of an amino acid containing amine group existing in natural RNA binding protein (Das and Frankel, Biopolymers, 70: 80-85, 2003). From the investigation of natural RNA binding proteins and binding peptides was confirmed that lysine or arginine which contains a large number of amine groups was included in the peptides and such proteins or peptides were already methylated considerably (Tan and Fankel, Proc Natl Acad Sci USA. 92, 5282-5286, 1995). It was additionally confirmed that RNA binding capacity was increased as methylation of arginine of RNA binding protein proceeded (Liu and Dreyfuss, Mol Cell Biol. 15, 2800-2808, 1995).\nThere have been a great numbers of reports on RNA binding capacity of a peptide containing amine group or RNA binding capacity depending on the methylation of amine group and natural RNA binding peptides and methylation of them. However, there was no report yet on synthesized RNA binding peptide or specific RNA binding capacity of a methylated peptide. Therefore, to obtain a peptide specifically binding to RNA, the present inventors prepared a peptide composed of 15 amino acids containing 7 alpha-helical lysines. In the meantime, to ensure the diversity of such peptides, a library was constructed by using the combination of methylated lysines. Then, the present inventors completed this invention by selecting peptides showing the strongest binding capacity to RRE-RNA of HIV-1 from those synthesized from the library. The peptide of the present invention thus has not only strong but also specific RRE RNA binding capacity, so that it can be used as a therapeutic agent for AIDS."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to an apparatus and a method of controlling a temperature of a matching material, and more particularly, to an apparatus and a method of maintaining a property of a matching material associated with an electromagnetic wave loss by controlling a temperature of the matching material used in an apparatus for generating a tomographic image. “This invention was funded by the MISP (Ministry of Science, ICT & Future Planning), Korea in the ICT R&D Program 2013 into an information and communication technology (ICT) research development program in 2013”\n2. Description of the Related Art\nIn recent times, research into microwave tomography (MT) using an electromagnetic wave is being conducted to develop a new technology for diagnosing breast cancer.\nAn apparatus for generating a tomographic image may immerse a breast in a predetermined matching material that allows an electromagnetic wave to readily penetrate the breast, and reconfigure a tomographic image of the breast based on wave scattering data that penetrates the breast.\nThe apparatus for generating the tomographic image may use amplitude and phase information of the electromagnetic wave that passes through the breast as key information to be used to reconfigure the tomographic image of the breast.\nHowever, the matching material has a property of causing a loss in the electromagnetic wave that passes through the matching material, and the property may vary based on a temperature of the matching material.\nWhen the property of the matching material is maintained, the amplitude and phase information of the electromagnetic wave that passes through the breast may be measured precisely based on the property of the matching material. However, when the property of the matching material changes, a time may needed to reconfigure the tomographic image or an error may occur in the reconfiguring.\nFurthermore, the property of the matching material may vary in response to the temperature of the matching material decreasing based on a period of time elapsing because the temperature of the matching material is higher than a general room temperature.\nAccordingly, there is a need for an apparatus and a method of preventing the amplitude and phase information of the electromagnetic wave that passes through the breast from being deteriorated due to the change in the property of the matching material based on a period of time elapsing."}
{"text": "Generally, a fistula is an abnormal connection or passageway between organs or vessels that normally do not connect. Fistulae can develop in various parts of the body. For example, types of fistulae, named for the areas of the body in which they occur, include anorectal fistula or fistula-in-ano or fecal fistula (between the rectum or other anorectal area and the skin surface), arteriovenous fistula or A-V fistula (between an artery and vein), biliary fistula (between the bile ducts to the skin surface, often caused by gallbladder surgery), cervical fistula (abnormal opening in the cervix), craniosinus fistula (between the intracranial space and a paranasal sinus), enteroenteral fistula (between two parts of the intestine), enterocutaneous fistula (between the intestine and the skin surface, namely from the duodenum or the jejunum or the ileum), enterovaginal fistula (between the intestine and the vagina), gastric fistula (between the stomach to the skin surface), metroperitoneal fistula (between the uterus and peritoneal cavity), perilymph fistula (a tear between the membranes between the middle and inner ears), pulmonary arteriovenous fistula (between an artery and vein of the lungs, resulting in shunting of blood), rectovaginal fistula (between the rectum and the vagina), umbilical fistula (between the umbilicus and gut), tracheoesophageal fistula (between the breathing and the feeding tubes) and vesicovaginal fistula (between the bladder and the vagina). Causes of fistulae include trauma, complications from medical treatment and disease.\nTreatment for fistulae varies depending on the cause and extent of the fistula, but generally involves surgical intervention. Various surgical procedures are commonly used, most commonly fistulotomy, placement of a seton (a cord that is passed through the path of the fistula to keep it open for draining), or an endorectal flap procedure (where healthy tissue is pulled over the internal side of the fistula to keep feces or other material from reinfecting the channel). Surgery for anorectal fistulae is not without side effects, including recurrence, reinfection, and incontinence.\nInflammatory bowel diseases, such as Crohn's disease and ulcerative colitis, are the leading causes of anorectal, enteroenteral, and enterocutaneous fistulae. The reported incidence of fistula in Crohn's disease ranges from 17% to 50%. Management of fistulae in patients with Crohn's disease continues to present an extremely challenging problem since many such fistulae do not respond to available treatments. Such fistulae and their recurrence are a very distressing complication that significantly reduces the quality of life of affected patients. Recent improvements in medical treatment (e.g., treatment with Infliximab) and expert surgical management have decreased the need for complicated surgery. However, many patients are not cured. Failure of fistulae to heal is probably due to the suboptimal quality of tissues that have been affected by Crohn's disease. Indeed, Crohn's fistulae provide a model system for wound healing under some of the worst possible conditions.\nAnother leading cause of fistulae is trauma, e.g. by rape, or by injuries sustained during childbirth, to the tissues of the vagina and the bladder and/or rectum leading to rectovaginal fistula and vesicovaginal fistula. Every year approximately 100,000 women across the developing world sustain such fistulae (also known as obstetric fistulae) during obstructed labor. During obstructed labor, the pressure of the baby's head against the mother's pelvis cuts off blood supply to delicate tissues in the region. The dead tissue falls away and the woman is left with a vesicovaginal fistula and sometimes a rectovaginal fistula. This hole results in permanent incontinence of urine and/or feces. The United Nations Population Fund (UNFPA) estimates the world's population of obstetric fistula sufferers at more than two million. This calculation could be a significant underestimate. Success rates for primary surgical repair range from 88 to 93 percent but decrease with successive attempts. Thus, a significant percentage of women have obstetrical fistulae that cannot be repaired surgically.\nNew therapies for fistulae are needed."}
{"text": "1. Field of the Invention\nThe present invention relates to a technology for a multiclient-support client-server type document management system.\n2. Description of the Related Art\nIn a document management system that is used by a plurality of client Personal Computers (PC) in a network, a document management server stores therein data to be displayed in the form of a list or the like on the display screen of each client PC. Some of the client PCs may refer to a list containing the same data in the document management server. In this case, a user of a client PC is not informed in real-time of changes in attributes or status of the-data caused by another user, and is not able to carry out optimal operation on the data in the list.\nFor example, suppose that the document management server includes multiple storages. If a user A transfers any data stored in a storage 1 to another storage 2, a user B who is referring to a list containing the data is not aware that the data has been transferred by the user A to the storage 2 because the list referred to by the user B indicates that the data is still stored in the storage 1. Due to this, the user B performs operation on the data as it is in the storage 1.\nIn a technology disclosed in Japanese Patent Laid-Open Publication No 2003-337813, a client-server type document management system includes a client PC and a document management server that are connected via a network. The document management server includes a notifying unit that allows, when management contents at a server side are changed by a client PC, a document manager to easily learn the change of the management contents. In a technology disclosed in Japanese Patent Laid-Open Publication No 2003-85024, a document management program and method are applied to a server connectable to a network for updatably managing a document file in the server on a World Wide Web (WWW) system.\nAs described above, in a document management system that is used by multiple client PCs on a network, a document management server stores therein data to be displayed in the form of a list or the like on the display screen of each client PC. Some of the client PCs may refer to a list containing the same data in the document management server. In this case, a user of a client PC is not informed in real-time of changes in attributes or status of the data caused by another user, and is not able to carry out optimal operation on the data in the list.\nOne approach to the problem is to send a message, from the document management server to each client PC, informing that a client PC has transferred data in the document management server. With the message, each client PC can update the list. However, there are document management systems in which the document management server does not issue such a message or the client PC is not able to process the message issued from the document management server. Thus, there is a need of a technology for allowing each client PC to voluntarily confirm status of data in the document management server."}
{"text": "The Kiosk was designed and developed to accommodate a need for a “stand alone unit” that houses an interactive computer monitor/touch screen display, for commercial trade shows and traveling exhibit applications. The requirements were to be lightweight, collapsible and shippable (UPS, FEDEX etc.) and yet maintain a “corporate” look. It was also required to have shelving for a CPU and CD/DVD player, keyboard, speakers for the audio aspect and a storage area for miscellaneous accessories. The kiosk also required that an operator be able to gain access to the equipment without completely disassembling the unit, so a locking door feature was incorporated into the present invention.\nFrom the foregoing, it will be appreciated that there is a need in the art to develop a sturdy lightweight foldable monitor stand with foldable shelves for easy portability and storage. The present invention is directed to overcoming one, or more, of the problem set forth above."}
{"text": "FEC (Forward Error Correction, forward error correction) is an error control technology in data transmission, and is widely used in communications field such as optical communication. The FEC has strong error correction capacity, and while a bit error rate before correction of long distance data transmission is greater than 2E-2 (representing that at most 1 bit error occurs when transmitting 2×102 bits), the bit error after FEC error correction can be kept below 1E-15.\nIn an optical communication system, a burst noise existing in a channel may lead to a burst bit error. The FEC has certain capacity of correcting the burst bit error; however, when the burst bit error caused by the burst noise in the channel exceeds the capacity of correcting the burst bit error of the FEC, the bit error may be spread, and a serious bit error rate after correction is generated.\nTo deal with the burst bit error, a channel interleaving solution is usually adopted in this field, that is, a continuous burst bit error is “scattered” to different FEC code words, so that the volume of the burst bit error in each code word is smaller than the capacity of correcting the burst bit error of the FEC. The specific method is: interleaving data after FEC coding at a transmitting end, correspondingly de-interleaving the interleaved data after being transmitted through a channel at a receiving end, and performing FEC decoding on the de-interleaved data. Three channel interleaving solutions exist in the prior art, namely, block interleaving, bit interleaved interleaving, and helical interleaving.\nFIG. 1A is a schematic diagram of a block interleaving solution. The block interleaving specifically adopts a “row-in and column-out” manner, that is, an information flow is buffered with an order of the row in a storage and then output in the column direction.\nFIG. 1B is a schematic diagram of a bit interleaving solution. The bit interleaved interleaving specifically mixes two information flows in a bit interleaved manner.\nFIG. 1C is a schematic diagram of a helical interleaving solution. The helical interleaving specifically adopts a “row-in and twill-out” manner, that is, an information flow is buffered with an order of the row in a storage and then output in the twill direction.\nHowever, when the foregoing three interleaving solutions match with a block convolutional decoder to correct the burst bit error, with increase of the interleave depth, the capacity of correcting the burst bit error of the foregoing three interleaving solutions may meet a bottleneck, and the implementation complexity of the three interleaving solutions is high."}
{"text": "Security is an important issue within the development of integrated circuits. Components and information within these integrated circuits require some way to provide access protection to protect the integrity of these systems. A security system of some form is usually implemented to prevent unauthorized access to particular locations of integrated circuits such as components (e.g. memory) of a data processing system. Trusted software developers are authorized users responsible for all applications on an integrated circuit such as a data processing system. They know of the existence, location, and use of every component and feature within the data processing system. In general, when a data processing system is released to the public, full disclosure of all the features is not given to the general public. They are only informed of those elements that apply to their needs and uses of the integrated circuit. Therefore, these authorized software developers, by having full access to the system, are responsible for all applications within the data processor while only limited information is given to the general users.\nA problem arises when unauthorized users access portions of the system that are not available to the general public. These unauthorized users can therefore achieve access to important components and violate the security of the system. This violation gives the unauthorized users important information concerning data or functions within the system, and this information can be detrimental to an integrated circuit such as a data processing system. A technique used to secure portions of a data processing system is to generate an exception fault anytime an access fault is detected. This exception fault is then handled by a software exception handler in a prior art manner. A large problem with a software exception handler arises when unauthorized users detect this software security feature and possibly trace the software in the handler. Unauthorized users can achieve unlimited access and can even override the handler in order to access protected components of the data processor or other integrated circuit. The software exception handler is very susceptible to penetration; therefore, a need exists for a more secure method to prevent unauthorized access to information stored in an integrated circuit."}
{"text": "1. Field of the Invention\nThis invention relates to and has among its objects the provision of novel apparatus for harvesting vegetable heads which are non-colinear with respect to the stalk line in a row of heads. Further objects of the invention will be evident from the following description wherein parts and percentages are by weight unless otherwise specified.\n2. Description of the Prior Art\nVegetable heads, such as cauliflower or cabbage, grown in the United States are presently harvested manually. The harvesters select the head to be harvested by visual observation according to size. The head is then manually cut and thrown into a suitable receptacle, e.g., a trailer drawn by a tractor.\nManual harvesting is fraught with disadvantages such as large consumption of time and expense, inevitable spoilage of mature vegetable heads, etc. For example, during peak growth periods a sufficient number of workers cannot be obtained to pick a cauliflower crop. As a result, the cauliflower heads, which can become over-mature one-day after they are ready for harvest, are lost as a salable product.\nMechanical means for harvesting lettuce heads are known and described in U.S. Pat. Nos. 3,300,954, and 3,300,955. In the known apparatus the maturity of the head is determined by compression rollers, one fixed and the other laterally movable. The compression rollers are composed of a plurality of cylindrical rollers mounted on shafts and projected downwardly. The compression rollers form a V-shaped configuration. A mature lettuce head will laterally displace the movable compression roller thus triggering a cutting mechanism which cuts the mature lettuce head. The cutting mechanism includes a pair of knife baskets vertically rotatable. Each knife basket is equipped with a knife and a basket or cradle portion. The knife severs the lettuce head and the cradle portion conveys the harvested head to a container.\nAlthough the known harvester works well for harvesting lettuce heads, it cannot be used for harvesting vegetable heads such as cauliflower and cabbage. A mature head of cauliflower, for example, is composed of a center or head of cauliflower engulfed in an abundance of large, loose-fitting leafy material. Thus, the cauliflower head is quite different from a lettuce head which is composed of leaves of lettuce tightly formed into a ball-like shape. Cauliflower and cabbage plants also have longer stems and grow further out of line from the row in which they are planted than do lettuce plants. That is to say certain of the heads ar non-colinear with the stalk line in a particular row. Generally, the stalks of the plant remain in a straight line corresponding to that formed when the seeds are planted. The plants, themselves, however, grow out of this stalk line, i.e., they are non-colinear therewith. Lettuce heads, on the other hand, will grow colinearly within the row.\nThe leaves surrounding the cauliflower and cabbage heads cause the known harvester to function improperly by fouling the compression rollers. In addition, the known selector cannot align itself with non-colinearly growing vegetable plants because one compression roller is fixed. Furthermore, the cutting mechanism is not adequate to sever mature heads of cauliflower or cabbage with their bulky leaf package."}
{"text": "1. Field of the Invention\nThis invention relates to a surface pressure distribution sensor which is suitable for detecting microscopic asperity patterns such as fingerprint patterns by using a flexible conductive film, and also relates to a manufacturing method of such a sensor.\n2. Description of the Related Art\nFIGS. 14A–14B show an example of an active matrix surface pressure distribution sensor for detecting fingerprint patterns. FIG. 14A is a plan view of the device, and FIG. 14B and FIG. 14C are cross sectional views taken along line D—D shown in FIG. 14A.\nA conventional sensor 200 for surface pressure distribution includes a substrate 201 which is made of a glass, a ceramic or the like, and a common electrode film 202. The device also has a number of TFTs (thin film transistors) 204a thereon as unit detection elements.\nEach of the unit detection elements 204 includes TFT 204a and contact electrode connected thereto. The unit detection elements 204 are arranged in the form of a matrix on the substrate 201. The active layers of the TFTs of the unit detection elements 204 are made of an amorphous silicon film. The contact electrodes 204b are made of ITO (indium tin oxide).\nThe common electrode film 202 is provided so as to face the substrate 201, and includes a flexible insulator film 202a and a conductive film 202b deposited on the rear side of the film 202a (TFT side). The common electrode film 202 is fixed on a sealing agent 203 applied around the substrate 201 so as not to be in contact with the substrate 201.\nAn example of a manufacturing method of this surface pressure distribution sensor will be described. After the TFTs are formed on the substrate 201, the sealing agent 203 made of a low temperature thermosetting resin is applied around the substrate 201 in order to affix the common electrode film 202 thereon. The common electrode film 202 is then affixed on the substrate 201 and subjected to a heat treatment. Consequently, the substrate 201 and the common electrode film 202 are fixed to each other.\nFIG. 14C shows an example of detecting fingerprint patterns by using this surface pressure distribution sensor. By placing a finger F to press slightly the top of the sensor 200, the common electrode film 202 as a whole is pressed down. However, the difference in pressure between the peaks and the valleys of the fingerprint pattern causes only the contact electrodes 204b of the unit detection elements 204 directly below and in the vicinity of the peaks to come into electrical contact with the common electrode film 202. On the other hand, the contact electrodes 204b of the unit detection elements 204 directly below and in the vicinity of the valleys of the fingerprint pattern are not in electrical contact with the common electrode film 202. Hence, the signals corresponding to the regions in which the common electrode film 202 and the unit detection elements 204 come into contact with each other are generated so as to detect fingerprint patterns.\nIt is known that a surface pressure distribution sensor with TFTs can be realized by the above-mentioned structure and manufacturing method. However, the reproducibility of such devices is poor when mass-produced."}
{"text": "To meet demands for smaller size and higher performance regarding high-frequency or radiofrequency (RF) component parts of cellular phones and the like, researches and developments of electric equipment using MEMS technologies are being vigorously carried out. In MEMS devices intended for radiofrequency uses, an RF switch, a variable capacitance or the like employs a movable portion that has a cantilever structure or a both-ends-supported beam (bridge shape) structure that is formed from a metal material structure of low resistance. The movable portion is displaced by piezoelectric drive, electrostatic drive, etc. For performing a desired function, the position of the movable portion should be stably controlled.\nAn MEMS switch is a mechanical switch that has a static or stationary electrode and a movable electrode facing each other and that performs on-off actions by driving or actuating the movable electrode to contact to or separate from the static electrode. The mechanical switch allows reduction in parasitic capacitance, and is low in loss, high in insulation and less in distortion in signal waveform, in comparison with switches that use semiconductor elements.\nIn a radiofrequency (RF) circuit, an MEMS capacitance is connected in series to or loaded on a RF line so as to define the frequency characteristic or adjust the distributed constant of the RF line. By using a variable capacitance, it is possible to change the resonance frequency or change the distributed constant. Generally, a variable-capacitance element has such a structure that a fixed electrode and a movable electrode face each other and the capacitance between the electrodes is changed by displacing the movable electrode.\nSuch a movable electrode is formed of a flexible metal structure that is formed by, for example, plating. Electrolytic plating requires an electricity feeding layer. For example, a stack of an adherence layer that provides adherence with a base and a seed layer made of the same material as a plate layer is formed by sputtering or the like. The adherence layer is formed of, for example, a metal layer of Ti, Cr, Mo, etc. The plate layer is formed, for example, of a highly electroconductive metal such as copper (Cu), gold (Au), etc.\nIn order to provide for a free space below the flexible metal structure, a method including forming a sacrificial film in a free space, forming a metal structure on top of the sacrificial film, and then removing the sacrificial film is employed. The sacrificial film that is used in this method may be, for example, a metal film of copper, aluminum, etc., an inorganic dielectric film of silicon oxide, silicon nitride, aluminum oxide, etc., or an organic dielectric film of a photosensitive resin, etc.\nFor example, in order to form a cantilever type movable electrode on a ceramics substrate, a pedestal portion is first formed by processes of forming an adherence layer/seed layer by sputtering or the like, forming a structure that defines a plating region by using a resist pattern or the like, forming a pedestal metal layer by electrolytic plating process, removal of a structure such as the resist pattern or the like, removal of the adherence layer/seed layer that is unnecessary, etc., and a sacrificial film that fills a free space is formed through formation of an adherence layer/seed layer by sputtering or the like, formation of a structure that defines a plating region through the use of a resist pattern or the like, formation of a sacrificial metal layer by an electrolytic plating process, removal of the structure of the resist pattern or the like, etc., and a cantilever structure is formed through formation of an adherence layer/seed layer by sputtering or the like, formation of a structure that defines a plating region through the use of a resist pattern or the like, formation of a movable beam portion metal layer by an electrolytic plating process, removal of the structure of the resist pattern or the like, etc. Thereafter, the sacrificial film and the unnecessary adherence layer/seed layer are removed, so that the cantilever type movable electrode is formed.\nThe cantilever structure, which is mainly formed from a good conductor such as gold or the like, includes an adherence layer and a seed layer as a base. A metal layer formed by sputtering and a metal layer formed by plating are different from each other in purity and the like and exhibit different physical properties even if the two layers are of the same metal. A stack of metal layers having different physical properties is a stack of metal layers whose thermal expansivities are different.\nAn electric equipment is subjected to a reflow step at about 260° C. or a temperature impact test at −20° C. to +80° C. Due to the different thermal expansivities, stress occurs between the stacked metal layers, and warpage or strain occurs. For example, a distal end of the cantilever structure warps and becomes displaced upward. In some cases, the distal end of the cantilever structure is displaced upward by 10 μm or more. As a result, the electric equipment fails to operate at a predetermined operating voltage. A cause of such warpage is considered to be that the structure is constructed of a stacked layer structure of layers of different metal materials. For example, the adherence layer and the plated layer are formed of different metal layers.\nFor example, the adherence layer may have a higher resistivity than the plated layer. If a contact surface of a switch is covered with an adherence layer, high resistance results. In order to reduce such resistance, it has been proposed to remove the adherence layer from the contact surface of a switch (e.g., Japanese Patent Application Publication No. 2007-196303 (JP 2007-196303 A) and Japanese Patent Application Publication No. 2009-252672 (JP 2009-252672 A))."}
{"text": "Pipe clamps are commonly used to join tubular components together; for example, pipes or tubular housings. These clamps can be used in a variety of applications with some clamps specifically designed for specific components or for use in specific applications, and others of a design intended to make them more generally or universally applicable. One such application of pipe clamps is in connecting pipes or other components in automotive exhaust systems. Often, these exhaust system applications require or at least desirably provide a joint between pipe ends that seals against exhaust gas leakage and that has good resistance against axial separation. One type of pipe clamp is a band clamp which is used with telescopically overlapping pipe ends, and another type is a pipe coupler which is used with end-to-end abutting pipe ends. Both types usually include a metal band to be placed and tightened over the pipe ends, and both types can include a sealing sleeve and/or a gasket to be sandwiched between the band and the pipe ends."}
{"text": "This application relates in general to separation techniques such as capillary electrophoresis and in particular, to apparatus and method for on-column derivatization in capillary electrophoresis.\nCapillary electrophoresis (CE) is rapidly emerging as one of the separation methods of choice in resolving a complex mixture into its constituents (see Jorgenson et al., Science 1983, 222, 266; Gordon et al., Science 1988, 242, 224; Ewing et al., Anal. Chem., 1989, 61, 292A). In this procedure an electric field is applied across a capillary structure with typical dimensions of 2-200 .mu.m inside diameter and 10-100 cm length. The medium is an electrolyte or a gel. Because the volumes are small, typically nanoliters of injected sample, a major challenge is to find suitable detection schemes. In past work, many detection schemes have been used, such as optical absorption (see Lauer et al., Anal. Chem., 1986, 58, 166), optical fluorescence (see Gassman et al., Science, 1985, 230, 813), electrochemical (see Wallingford et al., Anal. Chem., 1987, 59, 1762), conductimetric detection (see Huang et al., Anal. Chem.. 1987, 59, 2747), radioactivity (see Pentoney et al., Chromatog., 1989, 4809, 259, refractive index change (see Bruno et al., Appl. Spec., 1991, 45, 462, and mass spectrometry (see Olivares et al., Anal. Chem., 1987, 60, 1230). In each of these procedures, it may be advantageous to attach to the molecular constituent to be detected a label or tag that aids/enables its detection. Examples are labels for fluorescence, absorption, electrochemical detection and radioactivity.\nFrequently, the sample is derivatized off-column prior to being injected. In such pre-separation derivatization schemes, usually the sample to be separated is mixed with a labeling agent or compound in a vial where the compound or agent would react chemically with the sample to label or tag the sample. Such method of labeling or tagging is disadvantageous. Small samples will become diluted by the procedure since, for convenient handling, the vials used cannot be too small. After the labeling or tagging process has been completed, a portion of the labeled sample is then injected into a capillary column for separation. Thus, pre-column labeling requires two or more steps before the labeled sample is ready for separation. Post-column separation techniques have also been used, such as in U.S. Pat. No. 4,729,947 to Middendorf et al. Post-column labeling techniques share the same disadvantages as pre-column labeling techniques.\nIn view of the above-described drawbacks of off-column type derivatization techniques, on-column derivatization for fluorescence detection using a specially built T-shaped or \"cross\"-shaped structure in the capillary is proposed in U.S. patent application Ser. No. 07/235,953, filed Aug. 24, 1988 by Richard N. Zare et al., which is incorporated herein in its entirety by reference. See also the article entitled \"On-Line Connector for Microcolumns: Application to the On-Column o-Phthaldialdehyde Derivatization of Amino Acids Separated by Capillary Zone Electrophoresis,\" by Pentoney et al., Anal. Chem., 1988, 60:2625-2629.\nAnother type of on-column derivatization technique is disclosed in a poster presentation by U. R. Tjaden et al., entitled \"On-Line Derivatization and High Performance Capillary Electrophoresis,\" in the Thirteenth Symposium on Column Liquid Chromatography, Stockholm, Sweden, Jun. 25-30, 1989. In the technique proposed by Tjaden et al., a labeling reagent is added to the electrophoresis buffer in the capillary before the sample to be derivatized is introduced into the capillary column. Such technique also has a number of disadvantages. First, if fluorescent detection is used, the reagent present throughout the buffer in the capillary will cause background fluorescence which degrades the signal-to-noise ratio of the fluorescence detector. Furthermore, the reagent used must be carefully chosen so that the background fluorescence caused by the reagent will not be so high as to render fluorescence detection impossible. This severely limits the type of reagent that can be used and is undesirable.\nNone of the above-described off-column or on-column techniques is entirely satisfactory. It is therefore desirable to provide an improved on-column derivatization scheme in which the above-described difficulties are alleviated.\nIn particular, it is desirable to provide an improved on-column derivatization scheme where the unreacted reagent, dye or any other labeling compound would not reach the detector."}
{"text": "1. Field of Invention\nThe present disclosure relates to an electronic system and a routing method. More particularly, the present disclosure relates to a network system and a routing method.\n2. Description of Related Art\nWith advances in information technology, various kinds of network systems are widely used in our daily lives. Examples of such network systems include local area network systems, internet network systems, and data center network systems.\nTypically, a network system includes a plurality of nodes (e.g., switches). Upon reception of a packet, a node determines whether a characteristic of the received packet matches requirements of packet forwarding rules in the node itself, and forwards the received packet according to the packet forwarding rules. However, the packet forwarding rules may become invalid due to errors or due to having surpassed their effective dates. Therefore, when a node receives a packet, even if the characteristics of the packet are identical to a packet the node forwarded before, the node still has to plan a transmission path again to forward the received packet. As a result, the efficiency of the network system is decreased.\nThus, there is a need to redesign network systems to avoid unnecessary planning of transmission paths that would lead to decreasing the efficiency of the network systems."}
{"text": "The invention relates to sonic ranging and/or object detection systems, in general, and to such systems for aiding visually impaired persons, in particular.\nUltrasonic rangefinders for detecting the presence of or the distance to an object are well-known in the prior art. In, for example, U.S. Pat. No. 4,199,246 to MUGGLI, an ultrasonic rangefinder having a combination transmitting and receiving, capacitance-type, electrostatic transducer is incorporated in a photographic camera for the purpose of determining the distance to a subject and subsequently causing the adjustable focus lens of such a camera to be focused in accordance with a subject distance signal derived by said rangefinder.\nIn U.S. Pat. No. 4,280,204 to ELCHINGER, the disclosure of which is specifically incorporated herein, a conventional mobility cane for the blind incorporates ultrasonic object sensing apparatus in order to provide a blind, ambulatory cane user with the capability of remotely sensing the presence of movement-impeding obstacles within a spacial zone produced by said apparatus, whose size is infinitely variable. This sensing apparatus includes an adjustably mounted combination transmitting and receiving capacitance-type electrostatic transducer having an energy transmission pattern that approximates the size of said spacial zone.\nThe ultrasonic transducers described in the two above-mentioned patents transmit a directional, multiple-lobe pattern of ultrasonic energy whose contours are fairly well understood in the art. The multiple-lobe transducer pattern of a transducer with a circular backplate of 3.5 cm in diameter, for example, consists of a central lobe having a lobe angle of approximately 12.degree. at its half power point (-3dB) when operated at a frequency of 50 KHz, with said central lobe being generally symmetrical about a central axis and with a plurality of smaller magnitude side lobes that are also generally symmetrical about said central lobe axis. This electrostatic transducer, multiple-lobe pattern is described in much greater detail in an article by W. KUHL, et al., entitled \"Condenser Transmitters and Microphones with Solid Dielectric Airborn Ultrasonics\" in Acoustica, volume 4, 1954, pp. 519-532.\nLobe pattern shape of an electrostatic transducer of the type mentioned above is primarily a function of transducer operating frequency and transducer backplate diameter. The higher the operating frequency, the narrower is the central lobe angle and the lower the operating frequency, the wider is the central lobe angle. The central lobe angle of the object sensing apparatus disclosed in the above-mentioned ELCHINGER patent was purposely narrowed in order to provide a visually impaired person with a cue as to the direction of a particular detected object.\nA disadvantage associated with sonic sensing apparatus of the type described, for example, in said above-mentioned ELCHINGER patent is the need to regularly move such sensing apparatus from side-to-side in order to provide a mobility cane user with directional information with respect to laterally positioned objects. Constant movement of a cane of this type can become tiring to a cane user as well as inconvenient to manipulate. Failure to so manipulate such a cane can result in the failure to detect a laterally positioned object and subsequent injury to a blind mobility cane user walking into such an object.\nA primary object of the present invention is to provide apparatus for detecting the presence of frontally and laterally positioned objects that does not require the physical movement of said apparatus.\nAnother object of the present invention is to provide apparatus that can both determine the relative distance to frontally located objects and also detect the presence of laterally located objects.\nA further object of the present invention is to provide a mobility cane for the blind that will enable a blind user to determine the presence of and relative distance to frontally located objects as well as determine the presence of laterally located objects.\nOther objects, features and advantages of the present invention will be readily apparent from the following detailed description of the preferred embodiment thereof, taken in conjunction with the accompanying drawings."}
{"text": "Desire to conserve resources have led to efforts to incent users to carpool or use public transportation. Carpooling reduces pollution, and reduces the number of vehicles on the road. Aside from environmental benefits derived from carpooling, there are many “perks” and incentives associated with carpool. For example, special High Occupancy Vehicle (“HOV”) or carpool lanes have been instituted in various locations and countries to provide an incentive to carpool. Vehicles are permitted in the designated carpool lanes if they contain two or more passengers depending on the carpool lane requirement of a particular location. In addition, many public and private establishments offer special parking spots reserved for carpools, certain highways and bridges offer reduced tolls for carpoolers, and employees may subsidize vans and buses to incent employees to carpool. There have also been several “carpool challenge” contests were participants compete for prizes and fame by carpooling.\nRecognizing the significant benefits and practical incentives of carpooling many drivers will knowingly violate carpool regulations when they are driving alone. To deter the misuse of carpool incentives, violation of the carpool regulation may result in large fines. For example, in some locations a minimum HOV lane violation can subject a driver to a fine of nearly $500. In recent years, although there are many tools and mobile applications developed to help willing participants to organize carpooling and propose routes for carpools, there are no practical and easy management tools to help monitor and confirm the passengers' actual carpooling status other than by a visual confirmation of the passengers while they are in the vehicle. Therefore, there is a need for a user friendly and safe method and technology to confirm that individuals are actually carpooling in one single vehicle in real time so that carpool incentives can be properly distributed and administered, and the proper carpool participants can be awarded credit and incentives for their participation."}
{"text": "Concentrated oil-in water emulsions of liquid active ingredients or active ingredients dissolved in a solvent are commonly used in agricultural compositions due to certain advantages provided over other formulation types. Emulsions are water based, contain little or no solvent, allow mixtures of active ingredients to be combined into a single formulation and are compatible with a wide range of packaging material. However, there are also several disadvantages of such agricultural emulsions, namely that they are often complex formulations which require high amounts of surface-active agents for stabilization, are generally very viscous, have a tendency for Oswald ripening of the emulsion globules and separate over time. Therefore, improvements in such emulsion formulations are needed in the agricultural field.\nSeveral oil-in-water emulsion compositions for cosmetics and dermatological applications have been described in U.S. Pat. No. 5,658,575; U.S. Pat. No. 5,925,364; U.S. Pat. No. 5,753,241; U.S. Pat. No. 5,925,341; U.S. Pat. No. 6,066,328; U.S. Pat. No. 6,120,778; U.S. Pat. No. 6,126,948; U.S. Pat. No. 6,689,371; U.S. Pat. No. 6,419,946; U.S. Pat. No. 6,541,018; U.S. Pat. No. 6,335,022; U.S. Pat. No. 6,274,150; U.S. Pat. No. 6,375,960; U.S. Pat. No. 6,464,990; U.S. Pat. No. 6,413,527; U.S. Pat. No. 6,461,625; and U.S. Pat. No. 6,902,737; all of which are expressly incorporated herein by reference. However, although these types of emulsions have found advantageous use in personal care products, these types of emulsions have not been used previously with agriculturally active compounds, which are typically present in emulsions at much higher levels than cosmetic active ingredients.\nOne example of an agricultural oil-in-water emulsion composition that is suitable for agriculturally active ingredients that are liquid or soluble in suitable solvents at relevant storage temperatures is disclosed in U.S. patent application Ser. No. 11/495,228, the disclosure of which is expressly incorporated by reference herein.\nThe present invention is related to agricultural compositions comprising an oil-in-water emulsion, the oil-in-water emulsion composition having an oil phase and water phase, the oil-in-water emulsion composition comprising an oil adapted to form oily globules having a mean particle diameter of less than 800 nanometers, a polymeric modifier that is compatible with the oil phase, at least one agriculturally active compound, at least one non-ionic lipophilic surface-active agent, at least one non-ionic hydrophilic surface-active agent, at least one ionic surface-active agent, and water."}
{"text": "The invention relates to a method of severing or removing a biological structure, in particular bone, having a water-jet cutting system from which a severing medium under high pressure is discharged, and to a cutting-nozzle element and a water-jet cutting system.\nSuch methods are known on the market and are in use in many different forms and designs. In particular in medicine, it is known to sever, for example from outside, a bone by water-jet cutting. A disadvantage with this is that, in conventional water-jet cutting methods, the soft tissue, and not only the bone, is destroyed. The vascular system in the soft tissue at the bone is important in particular for the knitting of the bone or for the regeneration of the callus. It is therefore necessary during the water-jet removal or severing of biological substances, in particular of bones, to carry out the removal or severing of the bone as carefully as possible. In conventional water-jet cutting methods, the water is applied directly to the exposed bone via a cutting nozzle, in the course of which the vascular system in the bone is also damaged.\nAn arrangement for cutting by means of a liquid jet has been disclosed by EP 0 636 345 A1, in which arrangement an additional medium is added to a liquid jet by means of vacuum. In this case, pulsing of a liquid jet is produced in a handle, the liquid jet being discharged under pressure losses via an elongated cannula adjoining the handle."}
{"text": "During the last half century, considerable research has been devoted to the development of feedback engine control strategies which incorporate in-cylinder transducers for the measurement of selected values. In recent years, the primary focus of renewed interest has been in the possible development of practical on-board systems, for individual cylinder feedback trimming control of spark timing to MBT (Minimum Spark Advance For Best Torque) or to the knock limit, analogous to EGO (Exhaust Gas Oxygen) sensor feedback trimming control of air/fuel ratio to stoichiometry.\nFor the most part, prior art implementations have relied strictly on the use of vibration sensors in conjunction with elaborate signal processing to obtain the desired feedback result. See, for example, U.S. Pat. No. 5,040,510 entitled \"Method For Controlling Knocking In Internal Combustion Engines,\" issued to Krebs et al on Aug. 20, 1991 and assigned to Siemens Aktiengesellschaft. See also, U.S. Pat. Nos. 4,993,387 and 5,134,980 issued to Sakakibara et al on Feb. 19, 1991 and Aug. 4, 1992, respectively. Both the '387 and '980 patents are assigned to Nippondenso Co., Ltd. and relate to statistical based knock control systems for engines.\nAs previously indicated, each of the referenced prior art disclosures utilize vibration sensors to detect structure borne vibrations resulting from combustion chamber acoustic pressure oscillations (produced by knock or detonation). Vibration sensors of the type referenced above have generally proven inadequate as they also detect other structure borne vibrations and usually exhibit poor signal-to-noise ratios--particularly at high engine speeds. Similarly, the elaborate signal processing of the referenced prior art has proven expensive and temperamental and thus not desirable for non-laboratory based use."}
{"text": "A conventional saw blade clamping device 10 is disclosed in U.S. Pat. No. 5,647,133 and FIGS. 1 and 2, wherein the clamping device includes a body 11 having a slot 111 for receiving a blade 20 therein and a passage 112 defined in communication with the slot 111. A tube 14 has a rod 13 received therein and the rod 13 includes a rectangular end 131 which fits an inner periphery of the tube 14. A spring 15 is mounted to the rod 13 and biased between the rectangular end 131 and the inside of the tube 14. The rod 13 threadedly extends through the body 11 and further has a contact end 12 which is able to contact a side of the blade 20. The body 11 includes a serrated surface 113 and the tube 14 is firmly urged by the spring 15 to contact the serrated surface 113. A user has to pull the tube 14 to remove the tube 14 away from the serrated surface 113 and then rotate the tube 14 so as to move the rod 13 to urge the blade 20 or disengage from the blade 20. The user has to use a force that overcomes the force of the spring 15 and simultaneously, rotate the tube 14. This is inconvenient for the user to operate the tube 14 in two different directions. Furthermore, it is difficult to estimate the force that the contact end 12 contacts the blade 20.\nFIGS. 3 and 4 show the disclosure of U.S. Pat. No. 6,023,848 illustrating a blade clamping device 30 including a casing 31 and a base member 32 located in the casing 31. A first end of a biasing member 33 is connected to the base member 32. A blade 40 is engaged with the base member 32 and the biasing member 33 includes a protrusion portion 331 which urges against the blade 40. A lever 34 is pivotally connected to the casing 31 and includes an end that may push a free second end of the biasing member 33 to remove the protrusion portion 331 away from the blade 40. Although the biasing member 33 is easy to operate by operating the lever 34, the biasing member 33 quickly looses its biasing force after frequent operation by the lever 34."}
{"text": "In current display panels, a common pixel design is that a physical pixel is formed by three sub-pixels (a red sub-pixel, a green sub-pixel and a blue sub-pixel) or four sub-pixels (a red sub-pixel, a green sub-pixel, a blue sub-pixel and a white sub-pixel) for displaying, so that a physical resolution is a visual resolution. In practical applications, sometimes the visual resolution may be low, for example, in a process of continuously viewing images; and sometimes a high visual resolution is required, for example, at the time of viewing details of a fine image. Since the visual resolution of the display panel is fixed, requirements on different visual resolutions cannot be met."}
{"text": "1. Field of the Invention\nThe present invention relates to a technique for collectively connecting a plurality of grounding wires included in a wire harness for a vehicle to a given ground site inside the vehicle.\n2. Description of the Related Art\nHeretofore, as a ground connecting device for collectively connecting a plurality of grounding wires included in a wire harness for a vehicle, to a ground site of the vehicle, there has been known a type described in JP 10-208815A.\nFIG. 10 shows an outline of this device. The device comprises a harness-side connector 7 to be provided at a terminal end of a wire harness including a plurality of grounding wires, and a ground joint connector 1 to be fixed to a given ground site (in FIG. 10, a bolt 6) provided on a vehicle body 3. The harness-side connector 7 includes a plurality of non-illustrated female terminals to be attached to respective terminal ends of the grounding wires and a connector housing 8 for collectively holding the female terminals. The harness-side connector housing 8 has a plurality of built-in terminal locking portions for holding the female terminals respectively. The ground joint connector 1 includes a grounding conductor 5 and a connector housing 2 which holds the grounding conductor 5, the grounding conductor 5 integrally having a grounding terminal portion 4 to be fixed to the ground site and a plurality of non-illustrated male terminals provided inside the connector housing 2.\nAccording to this device, interconnecting the ground joint connector 1 and the harness-side connector 7 and fixing the grounding terminal portion 4 in the ground joint connector to the bolt 6 as the ground site establish a collective connection of the grounding wires to the ground site. Specifically, the female terminals held by the connector housing 8 of the harness-side connector 7 and the male terminals of the grounding conductor 5 held by the connector housing 2 of the ground joint connector 1 are fitted to each other respectively, thus electrically connecting the grounding wires to which the female terminals are attached to the ground site through the female terminals and the grounding conductor 5; simultaneously, the connector housing 8 of the harness-side connector 7 and the connector housing 2 of the ground joint connector 1 are fitted to each other, and this fitting is locked by engagement between respective engagement portions provided in the two connector housings 8 and 2, the lock keeping the female and male terminals fitted to each other respectively.\nHowever, this ground connecting device, occupying a large space, is difficult to use in a little space in a vehicle. Specifically, the harness-side connector 7 and the ground joint connector 1 of the device require the connector housings 8 and 2 for holding the terminals respectively; furthermore, the connector housings 8 and 2 occupy a large space as a whole for their mutual fitting and the lock of the fitting. To avoid interference between the connector housings 8 and 2 and the vehicle body 3, the connectors 7 and 1 are required to protrude in a large size from an inner surface of the vehicle body 3. Particularly, the case of connecting a grounding terminal 9 attached to an extra grounding wire W to the grounding terminal portion 4 so as to superimpose them to each other as shown in FIG. 10 requires a large gap size L between the vehicle body 3 and each of the connector housings 8 and 2 as shown in FIG. 10, in order to avoid the interference between the grounding terminal and each of the connector housings 8 and 2; this causes the entire device to occupy a larger space."}
{"text": "The invention relates to wound metallized film capacitors and more particularly to adding an additional auxiliary capacitor onto a first wound capacitor section to produce a dual value capacitor.\nOne method of making a dual value capacitor is disclosed in U.S. Pat. No. 3,921,041 issued to Stockman on Nov. 18, 1975. In this capacitor, one of the two basic metallized films has the metallized layer burned from its surface for a selected distance and an insulated sheet is inserted to encircle the capacitor at least once and to extend from one end of the roll. The winding is then completed around the insulated sheet and then both ends of the capacitor are shooped with one end providing the base plate and the other end, being divided into two sections by the extending insulated sheet, providing the other two sections of the dual value capacitor.\nThis method is not practical; however, when a small number of turns is desired to be added to the first section since there will be an insufficient end area on the auxiliary capacitor to attach a contact to. If a metal contact tab is added to the metallized film it will cause the film to burn away around the contact tab because of the small area involved. If the auxiliary section is end sprayed with metal, even if there is sufficient area to provide good electrical contact, the metallized film will burn away because of the small area involved.\nCurrent technology generally makes use of a separate capacitor if an auxiliary capacitor having a small number of turns is required in conjunction with a large capacitor. One such application of a large capacitor with an auxiliary small capacitor is in the fluorescent lighting ballast which requires a large 3.95 microfarad capacitor and a small 0.05 microfarad capacitor for a typical rapid start design. It would be desirable to utilize a wound metallized film construction to take advantage of the self-healing or clearing properties of the metallized film. When a short occurs the metallized film will self-heal and clear itself so that a short is only temporary and does not destroy the whole capacitor as happens in conventional metal foil capacitors. The need has thus developed for a wound metallized film capacitor having first and second sections which may be easily constructed to provide a compact capacitor having both a core capacitor and an auxiliary capacitor having a small number of turns within the same unitary structure."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus for applying a plurality of surgical fasteners to body tissue, and more particularly, to a surgical stapler including a staple cartridge assembly having a lockout mechanism for preventing refiring of the apparatus after the staples have been ejected from the cartridge.\n2. Description of the Related Art\nSurgical stapling apparatus for simultaneously applying a plurality of surgical fasteners to body tissue are well known in the art. Typically these apparatus include a fastener holder disposed on one side of the tissue to be fastened, an anvil assembly substantially parallel to the fastener holder on the other side of the tissue to be fastened, a mechanism for linearly translating the fastener holder and the anvil assembly toward one another so that the tissue is clamped therebetween, and a mechanism for driving the fasteners from the fastener holder so that the ends of the fasteners pass through the tissue and form finished fasteners as they make contact with the anvil assembly, thereby producing an array of finished fasteners in the tissue.\nIn common use are devices such as those disclosed in U.S. Pat. Nos. 4,354,628 and 4,665,916. More particularly, U.S. Pat. No. 4,354,628 discloses a surgical stapler apparatus for forming an array of surgical staples in body tissue including an anvil member against which the staplers are crimped, and a staple holder pivotally mounted adjacent one end of the anvil member.\nU.S. Pat. No. 4,665,916 discloses a surgical stapling apparatus comprising an anvil assembly against which fasteners are formed and a fastener holder pivotally mounted adjacent one end of the anvil assembly, a spacer member at the other end so constructed to displace tissue that would otherwise obstruct the spacer member from properly positioning the fastener holder relative to the anvil assembly to insure proper fastener formation, and a knife assembly to cut the tissue between the rows of formed fasteners.\nIn use, a surgeon selects the body tissue to be fastened, positions the instrument so that the tissue is between the anvil assembly and the fastener holder (or cartridge), then actuates the stapler. In some surgical applications, it is necessary to perform several stapling tasks and thus it is not uncommon for a surgeon to replace the cartridge several times during such procedures. In the course of an operation, however, a surgeon or nurse may, inadvertently try to reuse the apparatus with a spent cartridge in the apparatus or select a spent cartridge for placement into the apparatus. In such an instance, operation of the apparatus would be ineffective and would result in a prolongation of the procedure."}
{"text": "1. Field of the Invention\nThe subject invention relates to antennas. Particularly, the subject invention relates to microstrip antennas for circular polarization applications.\n2. Description of the Related Art\nAntennas for receiving signals from a satellite, such as Satellite Digital Audio Radio Service (SDARS) signals, are well known in the art. These antennas are routinely carried on vehicles for use with the vehicle's radio receiver. Typically, these antennas are mounted on a metallic roof of the vehicle such that the roof acts as a ground plane for the antenna. Furthermore, these antennas often have a bulky appearance that is not aesthetically pleasing from the outside of the vehicle.\nMany modern vehicles incorporate glass into their roof. The amount of glass used can range from a typical sunroof that provides glass over a small portion of the vehicle roof to a panoramic-style glass that spans the entire roof area of the vehicle. Unfortunately, the use of glass in vehicle roof structures reduces the amount of sheet metal that can be used as a ground plane for a satellite antenna. As such, typical satellite antennas suffer from lower performance when disposed on glass.\nTherefore, there remains an opportunity for an antenna that may be integrated with glass, such as a glass roof of a vehicle, for receiving signals from a satellite."}
{"text": "1. Field of the Invention\nThe present invention relates to a data readout circuit, a data readout method, and a data storage device.\n2. Description of the Related Art\nFIG. 1 is a circuit diagram showing the structure of a conventional ferroelectric memory. As shown in FIG. 1, the conventional ferroelectric memory comprises a word line WL, a bit line BL, a ferroelectric capacitor CF, n-channel MOS transistors 12, 14, 17A, 17B, 18A, 18B, and T1, p-channel MOS transistors T8 to T10, capacitors 19 and 22, and nodes NA and NB. The bit line BL has a bit line stray capacitance CBL.\nHere, the gate of the n-channel MOS transistor 14 is connected to the word line WL. One terminal of the source/drain of the n-channel MOS transistor 14 is connected to the bit line BL, while the other terminal is connected to the ferroelectric capacitor CF. A plate line CP is connected to one terminal of the ferroelectric capacitor CF.\nA voltage VCON is supplied to the n-channel MOS transistors 17A and 17B, while a reference voltage Vref is supplied to the source/drain of the n-channel MOS transistor 17A. A voltage VN is supplied to the gate of the n-channel MOS transistor 12, and a voltage RES is supplied to the gate of the n-channel MOS transistor T1. A voltage VP is supplied to the gate of the p-channel MOS transistor T8.\nIn the above ferroelectric memory, the single n-channel MOS transistor 14 and the single ferroelectric capacitor CF constitute one ferroelectric memory cell, as shown in FIG. 1. This ferroelectric capacitor CF holds digital information consisting of 1 or 0 in a non-volatile state by taking a reverse polarized state.\nNext, an operation to write the data in the ferroelectric memory cell will be described. When the information of xe2x80x9c1xe2x80x9d is written in the ferroelectric memory cell, the potential of the bit line BL serves as a ground potential. When the information of xe2x80x9c0xe2x80x9d is written in the ferroelectric memory cell, the potential of the bit line BL serves as a power potential Vcc. The word line WL is then activated, so that the n-channel MOS transistor 14 is energized and that the potential of the plate line CP changes from the ground potential to the power source potential Vcc and returns to the ground potential. Receiving the voltage, the ferroelectric capacitor CF shifts to a predetermined polarized state, and holds the information of xe2x80x9c1xe2x80x9d or xe2x80x9c0xe2x80x9d. When the data write operation is completed, the potential of the bit line BL is returned to the ground potential.\nNext, an operation to read out data from the ferroelectric memory cell will be described. In this case, the potential of the bit line BL serves as the ground potential. The word line WL is activated, so that the n-channel MOS transistor 14 is energized, and that the potential of the plate line CP shifts from the ground potential to the power source potential Vcc, thereby moving the charges polarized to the ferroelectric capacitor CF to the bit line BL. Here, the potential of the bit line BL greatly or slightly rises depending on the polarized state of the ferroelectric capacitor CF.\nFor instance, a latched sense amplifier circuit compares the potential of the bit line BL with the reference potential. In the initial state, the power source to the sense amplifier circuit is off, and when voltage is applied to the two input terminals, the power is supplied to the sense amplifier circuit. At this point, the input terminal having the potential higher than the other rises to the power source potential Vcc, and the input terminal having the potential lower than the other drops to the ground potential. By this sense amplifier circuit, data held by the ferroelectric capacitor CF can be read out.\nFIGS. 2A to 2I are timing charts showing the data read-out operation performed by the conventional ferroelectric memory shown in FIG. 1. As shown in FIGS. 2A and 2B, a voltage VCON and a signal RES are activated from 0 V (low level) to 3 V (high level) at time t1. By doing so, the potential of the bit line BL is initialized to 0 V, as shown in FIG. 2G. As shown in FIG. 2A, the voltage VCON is high until time t4.\nAs shown in FIG. 2C, the word line WL is activated at time t2, and the n-channel MOS transistor 14 is switched on. As shown in FIG. 2D, the potential of the plate line CP rises from 0 V to the power source potential (3 V) at time t3. Here, the potential of the bit line BL rises depending on the polarized charge amount of the ferroelectric capacitor CF, as shown in FIG. 2G.\nIn FIGS. 2G to 2I, each converted capacitance value, 0.2 pF, of the ferroelectric capacitor CF is indicated by a solid line, while each converted capacitance value, 0.05 pF, is indicated by a broken line. As can be seen from the timing charts, the larger the polarized charge amount, the higher the potential of the bit line BL. When the converted capacitance value of the ferroelectric capacitor CF is 0.2 pF, the potential of the bit line BL rises up to 0.5 V, which will be described more later.\nNext, as shown in FIG. 2E, the voltage VN to be supplied to the gate of the n-channel MOS transistor 12 that serves as a power switch for a sense amplifier is shifted to the high level at time t5. Here, as shown in FIGS. 2H and 2I, if the potential of the node NA (the potential of the bit line BL) is lower than the potential of the node NB (the reference voltage Vref), the potential of the node NA becomes 0 V while the potential of the node NB becomes equal to the reference voltage Vref, as indicated by the broken lines. On the other hand, the potential of the node NA is higher than the potential of the node NB, the potential of the node NA does not fluctuate, but the potential of the node NB shifts to 0 V, as indicated by the solid lines.\nNext, as shown in FIG. 2F, a voltage VP to be supplied to the gate of the p-channel MOS transistor T8 that serves as a VCC power switch for a sense amplifier is shifted to the low level at time t6. As shown in FIGS. 2H and 2I, if the potential of the node NA is lower than the potential of the node NB, the potential of the node NA is 0 V while the potential of the node NB becomes 3 V, as indicated by the broken lines. On the other hand, if the potential of the node NA is higher than the potential of the mode NB, the potential of the node NA becomes 3V while the potential of the node NB remains 0 V, as indicated by the solid lines.\nAs described above, after one of the potentials of the node NA and the node NB is shifted to 0 V while the other one of the potentials is shifted to 3 V, the potential of the node NA is transmitted via the bit line BL, so that the information stored in the ferroelectric memory cell is read out.\nIn the process of reading out information from the conventional ferroelectric memory shown in FIG. 1, the potential of the bit line rises depending on the polarization of the ferroelectric capacitor. Assuming that a cell capacitance value determined from the polarized charge amount of the ferroelectric capacitor and the voltage supplied between the electrodes is about 0.2 pF, the parasitic capacitance of the bit line is 1 pF, and the power source voltage is 3 V, the voltage of the bit line rises, by 0.5 V, which is calculated by 3 Vxc3x970.2 pF/(0.2 pF+1 pF), when the plate line rises from 0 V to 3 V. The potential of the bit line is shown in FIG. 2G. Accordingly, the voltage to be supplied to the ferroelectric capacitor becomes 2.5 V, which is calculated by 3 Vxe2x88x920.5 V.\nIntensive studies have been made on lowering the read-out voltage of the properties of the ferroelectric capacitor. However, polarized electric charge cannot adequately read out with a low read-out voltage, which results in inaccurate information read-out and reduction of read-out margins.\nMeanwhile, to reduce the power consumption of portable telephones and mobile electronic equipment, there has been a strong demand for lowering the power source voltage as well. If the bit line capacitance is increased, a rise of the voltage of the bit line can be reduced and a larger voltage difference can be applied across a ferroelectric capacitor. However, a read-out signal becomes smaller at the same time. As a result, in a latch-type sense amplifier circuit in a ferroelectric memory, a wrong operation is often caused due to an error in the input offset voltage.\nTherefore, it is necessary to design the latch-type sense amplifier circuit, so that a rise of the voltage of the bit line can be restricted to about 0.5 V. However, when the power source voltage drops to 2 V or 1 V, it is difficult to supply sufficient voltage across the ferroelectric capacitor.\nA general object of the present invention is to provide data storage devices in which the above disadvantages are eliminated.\nA more specific object of the present invention is to provide a data storage device that has high reliability and less power consumption, and a data read-out circuit and a data read-out method employed in the data storage device.\nThe above objects of the present invention are achieved by a data storage device that comprises: a memory cell connected between a plate line and a bit line; and a potential holding unit that maintains a potential of the bit line at a predetermined potential so as to prevent a fluctuation of the potential of the bit line even when a voltage is supplied to the plate line.\nWith this data storage device, the electric charge accumulated in the memory cell can be read out without fail.\nThe above objects of the present invention are also achieved by a data read-out circuit that reads out data from a memory cell, comprising: a charge accumulating unit that accumulates electric charge supplied; a charge transfer unit that transfers the electric charge accumulated in the memory cell to the charge accumulating unit in accordance with the data; and an amplifier unit that amplifies a voltage generated by the electric charge accumulated in the charge accumulating unit, and reads out the data from the memory cell.\nThe above objects of the present invention are also achieved by a data storage device that comprises: a bit line; a memory cell that is connected to the bit line; a charge accumulating unit that accumulates electric charge supplied; a charge transfer unit that transfers electric charge to the charge accumulating unit, the electric charge being accumulated in the memory cell based on stored data and then outputted onto the bit line at the time of reading out the data; and an amplifier unit that amplifies a voltage generated by the electric charge accumulated in the charge accumulating unit, and reads out data from the memory cell.\nWith this data storage device, data can be read out without fail, in accordance with the electric charge accumulated in the memory cell and outputted onto the bit line.\nThe above objects of the present invention are also achieved by a method of reading out data from a memory cell, comprising the steps of: transferring electric charge accumulated in the memory cell to a charge accumulating unit in accordance with the data; and amplifying a voltage generated by the electric charge accumulated in the charge accumulating unit, so as to reading the data from the memory cell.\nThe above and other objects and features of the present invention will become more apparent from the following description taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates to a powder accommodation container and an image forming apparatus."}
{"text": "1. Field\nThe present disclosure relates to an antiseptic applicator and method of use thereof, and more particularly, to a cap plug antiseptic applicator that requires the application of opposing forces to actuate release of a sealed solution, preferably an antimicrobial solution, from a self-contained reservoir toward a material arranged at a distal end of the applicator for receiving the solution.\n2. Description of Related Art\nAntiseptic applicators for the preparation of a patient prior to surgery, for example, are known and common in the prior art. Conventional applicators rely on various means of actuation to release a self-contained reservoir of antimicrobial solution for sterilization of the patient's skin. For example, a number of applicators are designed with a puncturing means. These applicators typically include a head with a spike, for example, and a sealed container or cartridge. A push or screw motion is employed to axially translate the head toward the sealed container so that the spike may pierce the sealed container and effectuate the release of the solution contained therein. Some examples of applicators using a puncturing means include U.S. Pat. Nos. 4,415,288; 4,498,796; 5,769,552; 6,488,665; and 7,201,525; and U.S. Pat. Pub. No. 2006/0039742.\nOther conventional applicators rely on breaking an internally situated frangible container or ampoule through the application of a one-way directional force or a localized application of pressure. The directional force is typically applied longitudinally to one end of the ampoule by a pushing motion designed to force the ampoule to break under a compressive stress, sometimes at a predetermined area of stress concentration. Alternatively, a pressure may be applied to a localized section of the ampoule through a squeezing motion designed to crush a section of the frangible ampoule in order to release the antimicrobial solution contained therein. Some examples of applicators using frangible ampoules in the manner discussed above include U.S. Pat. Nos. 3,757,782; 5,288,159; 5,308,180; 5,435,660; 5,445,462; 5,658,084; 5,772,346; 5,791,801; 5,927,884; 6,371,675; and 6,916,133.\nConventional antiseptic applicators, as described above, often require special packaging and/or handling during shipping and prior to use. For example, with the puncture type applicators, preventive measures are required to prevent an inadvertent push against either end of the device that may result in the puncturing of the sealed container and the premature discharge of the solution. A user must often use both hands to effectively overcome the preventive measures and activate the applicator for use. In addition, conventional antiseptic applicators often rely on the exertion of pressure on the walls of an applicator, for example, to break a frangible ampoule or squeeze the solution from the container toward an application material. The use of frangible ampoules requires special care to avoid breaking as a result of inadvertent pressure or dropping during shipping or prior to use. Furthermore, the components of a conventional applicator, such as the broken ampoule or the puncture spike, often impede the free flow of the solution from the container. There exists a need in the field for a novel antiseptic applicator that avoids the complications associated with conventional applicators, especially an applicator that will allow for effective one hand actuation and application of a solution without impediments to the free flow of the solution from the container to the application material."}
{"text": "The present invention relates generally to chemical mechanical polishing of substrates, and more particularly to methods and apparatus for monitoring a metal layer during chemical mechanical polishing.\nAn integrated circuit is typically formed on a substrate by the sequential deposition of conductive, semiconductive or insulative layers on a silicon wafer. One fabrication step involves depositing a filler layer over a non-planar surface, and planarizing the filler layer until the non-planar surface is exposed. For example, a conductive filler layer can be deposited on a patterned insulative layer to fill the trenches or holes in the insulative layer. The filler layer is then polished until the raised pattern of the insulative layer is exposed. After planarization, the portions of the conductive layer remaining between the raised pattern of the insulative layer form vias, plugs and lines that provide conductive paths between thin film circuits on the substrate. In addition, planarization is needed to planarize the substrate surface for photolithography.\nChemical mechanical polishing (CMP) is one accepted method of planarization. This planarization method typically requires that the substrate be mounted on a carrier or polishing head. The exposed surface of the substrate is placed against a rotating polishing disk pad or belt pad. The polishing pad can be either a “standard” pad or a fixed-abrasive pad. A standard pad has a durable roughened surface, whereas a fixed-abrasive pad has abrasive particles held in a containment media. The carrier head provides a controllable load on the substrate to push it against the polishing pad. A polishing slurry, including at least one chemically-reactive agent, and abrasive particles if a standard pad is used, is supplied to the surface of the polishing pad.\nOne problem in CMP is determining whether the polishing process is complete, i.e., whether a substrate layer has been planarized to a desired flatness or thickness, or when a desired amount of material has been removed. Overpolishing (removing too much) of a conductive layer or film leads to increased circuit resistance. On the other hand, under-polishing (removing too little) of a conductive layer leads to electrical shorting. Variations in the initial thickness of the substrate layer, the slurry composition, the polishing pad condition, the relative speed between the polishing pad and the substrate, and the load on the substrate can cause variations in the material removal rate. These variations cause variations in the time needed to reach the polishing endpoint. Therefore, the polishing endpoint cannot be determined merely as a function of polishing time.\nOne way to determine the polishing endpoint is to remove the substrate from the polishing surface and examine it. For example, the substrate can be transferred to a metrology station where the thickness of a substrate layer is measured, e.g., with a profilometer or a resistivity measurement. If the desired specifications are not met, the substrate is reloaded into the CMP apparatus for further processing. This is a time-consuming procedure that reduces the throughput of the CMP apparatus. Alternatively, the examination might reveal that an excessive amount of material has been removed, rendering the substrate unusable.\nMore recently, in-situ monitoring of the substrate has been performed, e.g., with optical or capacitance sensors, in order to detect the polishing endpoint. Other proposed endpoint detection techniques have involved measurements of friction, motor current, slurry chemistry, acoustics and conductivity. One detection technique that has been considered is to induce an eddy current in the metal layer and measure the change in the eddy current as the metal layer is removed."}
{"text": "Such methods, as well as corresponding devices are generally known. They are used for so-called cell harvesting, in other words, for separating and concentrating bacteria, for example, from liquid nutrient media. For this purpose, the centrifuges that are used are operated in the continuous flow mode. Solids are collected in the centrifuge rotor and the liquid components form the supernatant which is pumped out of the rotor and may be collected.\nIn this manner, very large quantities of suspension are centrifuged and separated. The liquid components, separated from the solids, are then separated into the high-molecular-weight and low-molecular-weight components in a separate process and/or a separate device using the so-called membrane filtration method. Various types of filters are used for the filtration technology, for example ultrafilters or microfilters which differ in the type of membrane used. Microporous membranes separate components in the 1/10 micron range, bacteria for example, while ultramembranes can separate much smaller components, proteins for example.\nThe known methods are also viewed as disadvantageous because they incorporate several separate method steps whose method parameters must be selected individually and therefore cannot be adjusted to one another, or can be adjusted only unsatisfactorily. In addition, these methods and/or the corresponding devices suffer from the disadvantage that the tanks with the components to be centrifuged or filtered must be moved around."}
{"text": "The present invention relates to an arrangement for mounting a restraint belt mounted vehicle seat to a vehicle floor and more particularly, to an arrangement for mounting a restraint belt mounted vehicle seat in an automotive vehicle that has an air bag system with multiple modes of inflation.\nA typical automotive vehicle seat has a seat bun or cushion that supports the buttocks and upper thigh region of a seated occupant. Adjustably and pivotally connected to the seat cushion is a seat back. The seat back supports the back region of a seated occupant. The seat cushion is connected to a seat riser. To allow the seat to have fore and aft adjustment, a seat adjuster is provided. The seat adjuster includes a seat channel, also referred to as a seat rail. The seat rail is slidably mounted on a lower rail, often referred to as a floor rail. The floor rail is typically connected to the floor pan of the vehicle. Typically, the seat rail is interlocked along its length with the floor rail to prevent vertical separation. To ease the sliding movement between the seat rail and the floor rail, ball bearings or rollers spaced between the seat rail and floor rail are provided. The seat rail has a spring-biased locking mechanism that engages with a connected or integral rack provided on the floor rail to lock the relative fore and aft position between the seat rail and floor rail. Typically, the vehicle seat will have two parallel sets of floor rails and seat rails. A master floor rail and seat rail combination will have a master latch which manipulates a slave latch unit on a parallel spaced slave floor and seat rail assembly.\nIn the most recent quarter-century, to facilitate vehicle safety, seatbelts have been added to vehicles. As a further enhancement of vehicle safety, three-point seatbelts have been provided which include shoulder restraints. Most front seatbelts have one end anchored to a B-pillar of the vehicle. The belt extends downward across the torso of a seated occupant through a loop. From the loop, the belt is routed across the seat occupant\"\"s lap and is then anchored to the vehicle floor. In a frontal crash, the load placed on the belt by a front seat occupant is mainly taken up by the B-pillar and/or the floor pan which the belt is anchored to.\nIn the most recent decade, a new type of anchoring system has been developed, commonly referred to as a belt-to-seat anchor restraint system. In the belt-to-seat anchor restraint system, one end of the belt is anchored to the B-pillar or to an upper region of the vehicle seat. The opposite extreme end of the belt is anchored to the upper portion of the seat riser which is fixably connected to the seat rail. The inner connection between the seat rail and the floor rail is strengthened to withstand the forces applied during a frontal crash situation.\nMany vehicles on the road today have airbags installed in steering wheels, dashboards, and more recently, doors. These airbags are designed to protect a vehicle occupant against both front and side impact collisions by rapidly inflating the airbag to absorb much of the collision energy that would otherwise be transferred to the occupant.\nSuch conventional airbags are inflated based on a single threshold test: if a predetermined vehicle deceleration occurs in a collision, airbag inflation is triggered. Thereafter, airbag deployment occurs at a predetermined inflation rate. Both the triggering threshold and the inflation rate are typically not modified based on the type of vehicle collision, or the many different occupant variables, such as occupant weight, occupant position at the moment of impact, etc.\nThere has been a desire to modify air bag deployment based upon occupant weight and position. Typically, the larger the occupant, the greater the desired inflative force. For smaller occupants, a lowered inflative force response is desired.\nThere are two major approaches to determine occupant weight on a vehicle seat. One approach is to have a pressure sensitive pad or bladder mounted somewhere within the seat cushion. Another approach is to place weight sensors between the floor pan and floor rail to sense the weight distribution on the vehicle seat and at a predetermined time sequence, inform the seatbelt inflater controller of the weight placed upon the vehicle seat. When using the weight sensor system, a new problem has occurred. The weight sensor typically adds a fixed link between the vehicle floor pan and the floor rail. Although this link is typically very strong in compression, there are limitations of this link in tension. Some vehicle occupant weight sensor systems rely upon a cantilevered support arrangement between the floor rail and floor pan of the vehicle. Cantilevered support arrangement weight sensors are typically very weak in tension.\nIt is desirable to provide a seat mounting system which allows for the utilization of a weight sensor element spaced between the floor rail of the vehicle seat and the floor pan of the vehicle, while at the same time allowing the vehicle to utilize a belt restraint seat anchoring system that anchors the restraint belt to the seat.\nIn a preferred embodiment, the present invention brings forth an arrangement for mounting a restraint belt mounted vehicle seat to a floor of an automotive vehicle. The arrangement includes a floor rail that supports the vehicle seat typically by a fore and aft seat adjuster that includes an interlocking seat rail. A load cell provides a supporting platform for the floor rail above the vehicle floor. Connected to the floor rail is a force transmittal member.\nIn a preferred embodiment of the present invention, the force transmittal member has an aperture. Inserted through the aperture is a headed fastener that is connected to the vehicle floor. Additionally, a lower riser is provided which is connected to the floor rail via the load cell. The lower riser is connected to and held in position on the vehicle floor by the aforementioned fastener. The head of the fastener restrains a capture member and connects the capture member to the vehicle floor. In a preferred embodiment the capture member has a generally U-shaped cross-section with extending flanges. The capture member is positioned adjacent the aperture in the force transmittal member.\nDuring normal operation, the compressive load of the vehicle seat is transmitted from the floor rail via the load cell to the lower riser and then to the vehicle floor. Upon a frontal crash situation, the floor rail will displace in an upward vertical motion and cause the force transmittal member to come into an interference situation with the capture member and thereby be retained to the vehicle floor. The floor rail will not be dependent upon the tensile strength of the load cell to return the vehicle seat to its position. The vehicle seat designer is now free to provide a load cell arrangement which can give sensory data to an air bag deployment system, while at the same time allow the vehicle seat to have a belt-to-seat mounting arrangement.\nIt is a feature of the present invention to provide an arrangement for mounting a restraint belt mounted vehicle seat to a floor of an automotive vehicle.\nIt is also a feature of the present invention to provide a method of retaining a vehicle seat to the floor of an automotive vehicle in a frontal crash situation wherein the vehicle seat utilizes load cells to sense the weight distribution upon the vehicle seat to inform a dual mode airbag inflation system.\nOther features of the invention will become more apparent to those skilled in the art upon a reading of the following detailed description and upon reference to the drawings."}
{"text": "The present invention relates to industrial fasteners and more particularly to an over-center toggle latch.\nLatches to join two members, for example, two panels, so that they may be separated and rejoined, have been developed over many years. Such latches are commonly used to secure container lids, trunk lids, panel doors and for industrial applications. A toggle latch provides the advantage of being relatively secure against accidental opening after it has been closed.\nAt the present time toggle latches are generally constructed with a bracket, an operating lever, a pivotable hasp (drawbar), an internal spring means if spring-loaded, and one or two pins pivotally mounting the lever and movable hasp on the bracket. That construction is relatively complex and expensive.\nA relatively simple fastener is shown in U.S. Pat. No. 4,049,301 to Peter Schenk which is entitled \"Toggle Latch\". The operating lever (handle) is pivotally mounted on the bracket by a first pin and the drawbar is mounted to the lever by a second pin. The drawbar is spring metal which is corrugated to provide spring action.\nIn U.S. Pat. No. 4,025,094 entitled \"Overcenter Latch\" a handle member is pivoted on one pin and a tension member pivots on a second pivot pin. In U.S. Pat. No. 3,026,133 to Swanson, the handle member pivots on a first pin, a coil spring in a housing pivots on a second pin on the handle member, and the movable hasp (link) pivots on a third pin on the handle member. Two patents (U.S. Pat. Nos. 3,847,423 and 4,243,255) to Rexnord, Inc. show the use of an internal coil spring in a toggle latch, both of which are relatively complex. U.S. Pat. No. 4,065,161, entitled \"Container Or Panel Clamp\", shows a simpler device in which the movable hasp (drawhook) is a wire member with a spring loop along its length. U.S. Pat. No. 2,820,995 to E. Schlueter, entitled \"Spring Loaded Lock Fastener\", shows a slidable hasp pivotally mounted, and spring-loaded, by spring wire pivot members. It is not a toggle latch, but rather operates by a cam mechanism."}
{"text": "This invention relates to an apparatus and method for aligning shafts. Particularly, this invention is useful in the new construction or maintenance fields in cases of setting and aligning the working or drive shafts of motors so that they can be coupled to the driven shafts of equipment, such as agitators, pumps, drive shafts, compressors and the like.\nWhenever a motor is attached to a piece of equipment to be operated the motor working or drive shaft must be attached to the equipment or driven shaft so that both are in perfect or near perfect alignment. If alignment is unsatisfactory, when the drive and driven shafts are coupled together by known conventional methods, the stresses resulting from operation in coupled misalignment will cause premature bearing failure in the motor or equipment or both. Often such failure will occur before operating speeds are attained in high speed electric motors or soon thereafter in lower speed motors. For these reasons engineering specifications often require alignment both horizontally and vertically of the drive and driven shafts within on thousandth (1/1000) of an inch.\nShaft alignment can be satisfactorily accomplished using a conventional alignment indicator which is fixed to one shaft and rides the other with a spring loaded stem capable of registering differences up to 200,000ths over two full turns on the indicator dial. When the shaft to which the indicator is attached is rotated through 360.degree., the variation in alignment is read at 90.degree., 180.degree., 270.degree. and 360.degree. or 0.degree.. A skilled man can then adjust one of the shafts until alignment indicator readings throughout the 360.degree. rotation are identical or within the specified tolerance. However, in order to have accurate readings, the indicator must be rigidly fixed to one of the shafts. Thus, one of the objects of the present invention is to provide an apparatus which can be rigidly fixed to a shaft to which the alignment indicator can be secured. Another object is to provide an alignment indicator clamp apparatus with more than one shaft engaging member. A further object of this invention is to provide an alignment indicator clamp apparatus which is easily attached and removed from the shaft without special tools. Another object of this invention is to provide an alignment indicator clamp apparatus which can handle shafts of widely varying diameter. A still further object of this invention is to provide an alignment indicator clamp apparatus which can be rigidly fixed to a shaft but is not cumbersome or bulky in size. These and other objects of the invention are accomplished by means of the apparatus of this invention."}
{"text": "1. Field of the Invention\nThe present invention pertains to typing recognition systems and methods, and more particularly to recognition of typing in air or on a relatively smooth surface that provides less tactile feedback than conventional mechanical keyboards.\n2. The Related Art\nTypists generally employ various combinations of two typing techniques: hunt and peck and touch typing. When hunting and pecking, the typist visually searches for the key center and strikes the key with the index or middle finger. When touch typing, the fingers initially rest on home row keys, each finger is responsible for striking a certain column of keys and the typist is discouraged from looking down at the keys. The contours and depression of mechanical keys provide strong tactile feedback that helps typists keep their fingers aligned with the key layout. The finger motions of touch typists are ballistic rather than guided by a slow visual search, making touch typing faster than hunt and peck. However, even skilled touch typists occasionally fall back on hunt and peck to find rarely-used punctuation or command keys at the periphery of the key layout.\nMany touchscreen devices display pop-up or soft keyboards meant to be activated by lightly tapping a displayed button or key symbol with a finger or stylus. Touch typing is considered impractical on such devices for several reasons: a shrunken key layout may have a key spacing too small for each finger to be aligned with its own key column, the smooth screen surface provides no tactile feedback of finger/key alignment as keys are struck, and most touchscreens cannot accurately report finger positions when touched by more than one finger at a time. Such temporal touch overlap often occurs when typing a quick burst of keys with both hands, holding the finger on modifier keys while striking normal keys, or attempting to rest the hands. Thus users of touchscreen key layouts have had to fall back on a slow, visual search for one key at a time.\nSince touchscreen and touch keyboard users are expected to visually aim for the center of each key, typing recognition software for touch surfaces can use one of two simple, nearly equivalent methods to decide which key is being touched. Like the present invention, these methods apply to devices that report touch coordinates interpolated over a fine grid of sensors rather than devices that place a single large sensor under the center of each key. In the first method, described in U.S. patent application Ser. No. 09/236,513 by Westerman and Elias, the system computes for each key the distance from key center to the sensed touch location. The software then selects the key nearest the finger touch. In the second method, described in U.S. Pat. No. 5,463,388 to Boie et al., the software establishes a rectangle or bounding box around each key and decides which, if any, bounding box the reported touch coordinates lie within. The former method requires less computation, and the latter method allows simpler control over individual key shape and guard bands between keys, but both methods essentially report the key nearest to the finger touch, independent of past touches. Hence we refer to them as ‘nearest key’ recognizers.\nUnlike touchscreens, the multi-touch surface (MTS) described by Westerman and Elias in Ser. No. 09/236,513 can handle resting hands and temporal finger overlap during quick typing bursts. Since the MTS sensing technology is fully scalable, an MTS can easily be built large enough for a full-size QWERTY key layout. The only remaining barrier to fast touch typing on an MTS is the lack of tactile feedback. While it is possible to add either textures or compressibility to an MTS to enhance tactile feedback, there are two good reasons to keep the surface firm and smooth. First, any textures added to the surface to indicate key centers can potentially interfere with smooth sliding across the surface during multi-finger pointing and dragging operations. Second, the MTS proximity sensors actually allow zero-force typing by sensing the presence of a fingertip on the surface whether or not the finger applies noticeable downward pressure to the surface. Zero-force typing reduces the strain on finger muscles and tendons as each key is touched.\nWithout rich tactile feedback, the hands and individual fingers of an MTS touch typist tend to drift out of perfect alignment with the keys. Typists can limit the hand drift by anchoring their palms in home position on the surface, but many keystrokes will still be slightly off center due to drift and reach errors by individual fingers. Such hand drift and erroneous finger placements wreak havoc with the simple ‘nearest key’ recognizers disclosed in the related touchscreen and touch keyboard art. For example, if the hand alignment with respect to the key layout drifts by half a key-spacing (˜9 mm or ⅜″), all keystrokes may land half-way between adjacent keys. A ‘nearest key’ recognizer is left to choose one of the two adjacent keys essentially at random, recognizing only 50% of the keystrokes correctly. A spelling model integrated into the recognizer can help assuming the typist intended to enter a dictionary word, but then actually hinders entry of other strings. Thus there exists a need in the touchscreen and touch keyboard art for typing recognition methods that are less sensitive to the hand drift and finger placement errors that occur without strong tactile feedback from key centers.\nFor many years, speech, handwriting, and optical character recognition systems have employed spelling or language models to help guess users' intended words when speech, handwriting, or other input is ambiguous. For example, in U.S. Pat. No. 5,812,698 Platt et al. teach a handwriting recognizer that analyzes pen strokes to create a list of probable character strings and then invokes a Markov language model and spelling dictionary to pick the most common English word from that list of potential strings. However, such systems have a major weakness. They assume all user input will be a word contained in their spelling or language model, actually impeding entry of words not anticipated by the model. Even if the user intentionally and unambiguously enters a random character string or foreign word not found in the system vocabulary, the system tries to interpret that input as one of its vocabulary words. The typical solution is to provide the user an alternative (often comparatively clumsy) process with which to enter or select strings outside the system vocabulary. For example, U.S. Pat. No. 5,818,437 to Grover et al. teaches use of a dictionary and vocabulary models to disambiguate text entered on a ‘reduced’ keyboard such as a telephone keypad that assigns multiple characters to each physical key. In cases that the most common dictionary word matching an input key sequence is not the desired word, users must select from a list of alternate strings. Likewise, users of speech recognition system typically fall back on a keyboard to enter words missing from the system's vocabulary.\nUnfortunately, heavy reliance on spelling models and alternative entry processes is simply impractical for a general-purpose typing recognizer. Typing, after all, is the fallback entry process for many handwriting and speech recognition systems, and the only fallback conceivable for typing is a slower, clumsier typing mode. Likewise, personal computer users have to type into a wide variety of applications requiring strange character strings like passwords, filenames, abbreviated commands, and programming variable names. To avoid annoying the user with frequent corrections or dictionary additions, spelling model influence must be weak enough that strings missing from it will always be accepted when typed at moderate speed with reasonable care. Thus a general-purpose typing recognizer should only rely on spelling models as a last resort, when all possible measurements of the actual typing are ambiguous."}
{"text": "A particular type of inflatable vehicle occupant protection device is commonly referred to as an air bag. An air bag is stored in a vehicle in a folded, uninflated condition at a location adjacent to the vehicle occupant compartment. When the air bag is to be inflated, inflation fluid is directed to flow from a source of inflation fluid into the air bag. The inflation fluid inflates the air bag from the folded, uninflated condition to an unfolded, inflated condition in which the air bag extends into the vehicle occupant compartment.\nWhen the air bag is being inflated into the vehicle occupant compartment, it engages an occupant of the vehicle to help protect the occupant from a forceful impact with parts of the vehicle. The manner in which the air bag engages the occupant is determined in part by the configuration and location of the air bag relative to the occupant as the air bag unfolds and moves into the vehicle occupant compartment."}
{"text": "1. Field of the Invention\nThe present invention relates to a power driven screwdriver having a clutch mechanism for transmitting rotation of a drive motor to a spindle with a driver bit.\n2. Description of the Prior Art\nIn a power driven screwdriver, a clutch mechanism is provided for transmitting and disconnecting the rotation of a drive motor to a spindle with a driver bit. The clutch mechanism is normally constructed as a claw clutch and includes a pair of clutch members, one of which is mounted on the spindle and the other of which is mounted on a main gear driven by the drive motor. The spindle is movable in an axial direction for engaging and disengaging the clutch members. With such a clutch mechanism constructed by a simple claw clutch, since the rotation of the spindle is restrained, for example, at the completion of a screw driving operation, the clutch mechanism temporarily repeats its engaging and disengaging operation. This will generate clanging sounds, giving unpleasant feeling to the operator, and cause early wear of the clutch mechanism.\nU.S. Pat. No. 4,655,103 discloses a power driven screwdriver including stopper for adjusting the driving\na amount of a screw by a driver bit. A claw clutch mechanism is provided between a driver shaft and a spindle movable in an axial direction. The claw clutch mechanism includes a first and a second clutch member formed on the driver shaft and the spindle, respectively. A clutch disc is interposed between the driver shaft and the spindle and includes a third and a fourth clutch member for engagement with the first and second clutch members respectively. A spring is interposed between the first and third clutch members for normally keeping them at a disengaging position. The second and fourth clutch members includes relief portions which serves not to transmit rotation. When the stopper abuts on a work to be screwed, the driver shaft continues rotation while the rotation of the spindle is prevented. This may cause the operation of the relief portions of the second and fourth clutch members to positively disengage the first and the third clutch members with the aid of the spring.\nU.S. Pat. No. 4,809,572 discloses a power driven screwdriver including a stopper sleeve for adjusting the driving amount of a screw and a claw clutch mechanism having a pair of clutch members, one of which is mounted on a main gear driven by a drive motor, while the other of which is mounted on a spindle. A spring is provided for normally keeping the clutch member of the spindle out of engagement with the clutch member of the main gear. A control mechanism is provided between the spindle and the clutch member mounted on the spindle. The control mechanism includes oblique recesses and a ball for engagement with the recesses. With such construction, when the stopper sleeve abuts on a work to be screwed, the main gear continues its rotation while the rotation of the spindle is prevented. In this stage, the control mechanism operates to positively move the clutch member of the spindle out of engagement with the clutch member of the main gear with the aid of the spring.\nHowever, with the above prior U.S. Patents, the operation of the clutch mechanism must accompany a reciprocal movement of the spindle at a long distance. In general, a power driven screwdriver is provided with a seal member for sealing between a spindle and a housing to prevent entry of dust within the housing. In case the spindle reciprocally moves at a long distance, the dust may be absorbed into the housing through the ga between the seal member and the spindle or the housing by the pumping effect. Thus, when the spindle moves into the housing, negative pressure will be created in the housing. Such dust entered into the housing may cause early wear or damage of the clutch mechanism or bearings disposed within the housing.\nFurther, with the clutch mechanism of the above U.S. Patents, after the stopper or the stopper sleeve has abutted on the work, no further driving operation cannot be made even if the driving of a screw wa insufficient."}
{"text": "Such a cartridge is known from DE 10 2008 057 443 A1, where the functional element is a valve device.\nIt is known from U.S. Pat. No. 4,391,590 B1 that a functional element is designed as a cap, which is pulled over a cannula opening of a cannula duct to close the cannula section. It is disadvantageous here that the cartridge and the cap for the cartridge must be manufactured in two mutually independent manufacturing steps, wherein the cap is placed, as a rule, by hand by a person or in an automated manner on the cannula opening of the cannula section. This causes high manufacturing costs. In addition, with the cap already placed, there is a risk of air inclusions during the filling of the cartridge with the dental material, as result of which the shelf life of the dental material may be reduced and/or the quantity of the filling may show undesired variations in a comparison of a plurality of cartridges. Such air inclusions may lead to the loss of the cap, especially during transportation, because of the expansion of the air, as a result of which the storage stability is reduced. Even though the inclusion of air can be reduced when filling the cartridge without cap, there is a risk now that dental material will escape from the cannula opening of the cannula section, as a result of which dental material will be lost. This leads to higher manufacturing costs. Such cartridges are intended for single-time use especially in the field of dentistry. However, there is a risk when using caps that the cannula section will be reclosed with the cap in order to use a residual material that is preset later. There is a risk of contamination of the dental material and/or of an increased risk for infection because of its undesired reclosing of the cartridge that was once opened.\nA cartridge, in which fibers or a flocking are connected to the cartridge in the area of an outlet of the cannula section, is known from U.S. Pat. No. 6,059,570. This functional element is used to apply, spread and/or burnish the dental material. It is disadvantageous here that the application of the fibers or of the flocking is carried out in an independent manufacturing step and fully independently from the manufacture of the cartridge. It is also disadvantageous that the cannula section is rigid in the area of the fibers or flocking. As a result, there is a risk that a treatment with the functional element is perceived by a patient as being unpleasant and/or painful. In addition, there is a risk that undesired injuries will develop because of the rigid design. The spreading and/or burnishing of the dental material is also made difficult by the rigid design of the cartridge and of the cannula section.\nFurthermore, it is disadvantageous in prior-art cartridges that there is a risk of an especially abrupt rupture of material in case of an overstressing due to an excessively strong force or an excessively high pressure being applied to press the dental material out of the cartridge and/or the reservoir."}
{"text": "A polymer electrolyte fuel cell in which a hydrogen-containing fuel gas and oxygen-containing oxidizing gas are supplied to an anode and cathode, respectively, and an electromotive force is generated by an electrochemical reaction occurring at both poles is generally constituted by sequentially laminating a bipolar plate, a gas diffusion electrode substrate, a catalyst layer, an electrolyte membrane, a catalyst layer, a gas diffusion electrode substrate, and a bipolar plate. The gas diffusion electrode substrate is required to have high gas diffusivity for allowing a gas supplied from the bipolar plate to be diffused into the catalyst layer and high water removal performance for discharging liquid water generated by the electrochemical reaction to the bipolar plate, as well as high electrical conductivity for extracting generated electric current, and gas diffusion electrode substrates composed of carbon fibers and the like are widely used.\nHowever, the following problems are known: (1) when the polymer electrolyte fuel cell is operated at a relatively low temperature of below 70° C. in a high current density region, as a result of blockage of the electrode substrate by liquid water generated in a large amount and shortage in the fuel gas supply, the fuel cell performance is impaired (this problem is hereinafter referred to as “flooding”) ; and (2) when the polymer electrolyte fuel cell is operated at a relatively high temperature of 80° C. or higher, as a result of drying of the electrolyte membrane due to water vapor diffusion and a reduction in the protonic conductivity, the fuel cell performance is impaired (this problem is hereinafter referred to as “dry-out”). In order to solve these problems of (1) to (2), various efforts have been made. A method of improving gas diffusivity and water removal performance by forming a microporous part on the surface of the gas diffusion electrode substrate, and forming pores in the microporous part is the basic solution to these problems.\nPatent Document 1 discloses that stable fuel cell performance can be obtained in a low humidity condition and high humidity condition by having a structure in which a carbon porous material, i.e., a microporous part, is impregnated in an electrode substrate, and the density of the impregnated layer is set to a predetermined range. However, by the structure in which a microporous part is impregnated in an electrode substrate, obtained by the above method, high gas diffusivity and high water removal performance cannot be simultaneously satisfied, and particularly, fuel cell performance has been insufficient at low temperatures.\nPatent Document 2 discloses a technology to form a through hole by putting a large quantity of pore-forming particles into the inside of the microporous part, and obtaining high performance in the drying conditions and humidified conditions by separating the paths of water and gas. However, while water removal performance is improved by the microporous part in the method disclosed in Patent Document 2, there is a problem that discharged water accumulates in carbon paper and inhibits diffusion of gas, and sufficient properties could not be obtained.\nAs described above, a variety of efforts have been made; however, one that can be satisfied as a gas diffusion electrode substrate which has excellent anti-flooding characteristic particularly at low temperatures without deteriorating anti-dry-out characteristic is yet to be discovered."}
{"text": "Typical electrical connector assemblies include a pair of connector housings of dielectric material, such as male and female connector housings, which mount complementarily interengaging terminals so that the terminals complete electrical circuits when the housings are fully mated in some applications, the male and female connector housings are relatively easy to separate or unmate to disconnect the electrical circuit. In other applications, it may be desirable to prevent separation or unmating of the connectors if at all possible.\nFor instance, in automotive applications, it may be desirable to unmate a pair of electrical connectors for purposes of servicing a vehicle, but unintentional separation of the connectors could lead to serious consequences. Similarly, in various appliance applications, such as refrigerators or food freezers, again it may be desirable to separate or unmate the connectors for purposes of servicing the appliances, but unintentional separation of the connectors could lead to loss or spoilage of the appliance contents.\nConsequently, various approaches have been made to design mating electrical connectors with security locking systems which prevent separation of the connector housings of a male and female connector assembly.\nFor example, many electrical connector assemblies employ latching devices between the housings of the male and female connectors. The latching devices often are provided in the form of flexible latch arms on one connector housing for snapping into engagement with latch bosses or detents on the mating connector housing. It has been found that such flexible latch arms, themselves, can be unintentionally separated through handling or by other means. Consequently, security locking keys have been employed in conjunction with such latching devices to prevent the latching devices, themselves, from becoming disengaged.\nHeretofore, one of the problems with security locking key systems has been that the locking keys are difficult, if not impossible, to remove should it be desirable to separate the connectors, without destroying the keys or associated portions of the connector housings. It would be desirable to provide a simple security locking key system wherein the locking keys can be easily removed to separate the connectors if desired.\nAnother problem with security locking key systems in electrical connector applications is that the keys often are misplaced or simply difficult to handle while manipulating the connector housings of the connector assembly.\nThis invention is directed to providing a simple, yet effective security locking key system for electrical connectors, wherein a locking key can be readily removed without destroying the key or associated portions of the electrical connectors, and also to providing a security locking key system wherein the locking key is tethered to the electrical connector assembly so that it cannot be misplaced and is readily available for use."}
{"text": "This invention relates to an infrared detecting element and also an infrared imaging device.\nSome infrared detectors use Si crystals and detect infrared rays having wavelengths equal to or longer than several micrometers. Such infrared detectors are of two types, the first type being produced by doping impurities into the Si crystals and the second type using heterojunction barriers.\nInfrared Detectors II, Chapter 2, Semiconductors and Semimetals, written by P. R. Bratt, published from Academic Press in 1977, discloses the first-type infrared detectors.\nJapanese published unexamined patent application 61-241985 discloses the second-type infrared detector. The documents \"3P79\" of the lecture in the thirty-third spring meeting of Applied Physical Society of Japan in 1986 also discloses the second-type infrared detector.\nThese two types of infrared detectors are useful for infrared two-dimensional imaging devices of a monolithic type. The first-type infrared detectors have the following drawback. Since the quantity of doped impurities is limited, the detector sensitivity is low and the detected wavelength is fixed in dependence on the type of the impurities. Accordingly, it is impossible to maximize the detector sensitivity at an arbitrary wavelength. The second-type infrared detectors are free from such a drawback."}
{"text": "1. Field of the Invention\nThe present invention relates generally to modern control systems and, more particularly, to negative feedback loops in such systems.\n2. Description of the Art\nFIG. 1 illustrates a known control system utilizing a negative feedback loop in a low drop-out (LDO) amplifier application 100. This particular application 100 is configured as an LDO regulator circuit. An LDO regulator is a circuit that provides a well-specified and stable DC voltage. The lowest value of differential (input/output) voltage at which the control loop stops regulating is called the dropout voltage. Modern applications such as communication electronics and other battery-powered portable devices require a low dropout voltage and low quiescent currents for increased power efficiency. LDO regulators meet both of these design needs.\nAt the input stage, a reference input signal VREF is fed into the inverting input of a dual stage amplifier 104. The output from the amplifier controls a field effect transistor (FET) Q1 that acts as a switch for supplying current from the power source VDD to the load (modeled as a resistor RL in the figure). Some of the current flowing between the source and the drain of Q1 is then fed back through a simple RC filter network into the non-inverting input of the amplifier 104. This feedback signal is called VFB. The RC filter network comprises capacitor C1 and resistors R1 and R2. C1 AC-couples the output back into amplifier 104. Resistors R1 and R2 are configured in a voltage divider with R2 connected to ground. The ratio between the values of R1 and R2 may be adjusted to set the output voltage, VOUT, to a desired value.\nVOUT is fed back through the RC filtering network yielding signal VFB at the non-inverting input of the amplifier. Typically, differential amplifiers are used in modern electronic circuits. Differential amplifiers amplify the voltage difference between two input signals. When the output of a differential amplifier is connected to its inverting input and a reference voltage signal is applied to the non-inverting input, the output voltage of the op-amp closely follows that reference voltage. As the amplifier output increases, that output voltage is fed back to the inverting input, thereby acting to decrease the voltage differential between the inputs. When the input differential is reduced, the amplifier output and the system gain are also reduced. In FIG. 1, because amplifier 104 is a dual-stage amplifier, the reference signal is shown connected to the inverting input rather than the non-inverting input. Nevertheless, because the output is fed back in a manner that reduces the system gain, the result is negative feedback, sometimes called degenerative feedback.\nNegative feedback is often employed to stabilize a control system when the system exhibits a gain from the input to the output. The output stage 120 in this LDO application is modeled by load resistor RL and an output capacitor C0 which is needed to deliver an instantaneous current to a dynamic load. C0 has a characteristic equivalent series resistance (ESR) modeled by a series resistor RESR. ESR is an effective resistance that is used to describe the resistive part of the impedance of certain electrical components such as capacitors.\nAn important characteristic of this type of control circuit is the ratio between the output and input signal amplitudes, known as the transfer function. The transfer function for any given system is used to model the gain of the system as a function of the input signal frequency. Such control systems are often designed to meet the specifications of a transfer function. The frequency response of the control system is completely described by its transfer function. As such, the stability of a system over a range of input signal frequencies may be predicted based upon properties of its transfer function known as poles and zeros. A pole is a root of the polynomial denominator of a transfer function; a zero is a root of the polynomial numerator.\nIn designing stable systems, one important consideration is the shift in phase that a signal undergoes as it passes through the system. Poles and zeros are associated with these shifts in phase. If the signal accumulates a shift in phase of 180 degrees, the shift causes the negative feedback to become positive feedback. This is problematic when the system is operating at greater than unity gain as positive feedback will drive the system to an unstable oscillatory state. In order to maintain the stability of the control system, designers often build in a phase shift buffer, called a phase margin. For example, a 50 degree phase margin ensures that the signal never undergoes a phase shift of more than about 130 degrees (i.e. it never comes within approximately 50 degrees of a 180 degree phase shift). 50 degrees is a typical value of a phase margin in an LDO design; however, a 50 degree phase margin is not a requirement for stability and smaller phase margins of 45 degrees or lower may suffice. Furthermore, although a design goal may be to maintain a particular phase margin, the actual performance of a system may be less than the nominal phase margin value. The nominal value of the phase margin is chosen to meet the specifications of a particular design and may vary significantly.\nBoth poles and zeros can be introduced into the transfer function describing the control loop by inserting various electronic components into the loop. For example, a dual-stage amplifier will create two poles in the transfer function. The addition of poles and zeros into the frequency response of a system must be taken into account in order to design a system with a bounded (finite) output. Unwanted or unavoidable poles and zeros can create significant challenges when trying to stabilize a control system over a range of operating frequencies.\nPreviously, efforts have been made to stabilize a control system by designing the system so that troublesome poles only affect the system negligibly over the operating frequency range. This approach limits the designer to specific component values and configurations. For example, an output stage may include a capacitor having an ESR which adds a zero to the transfer function at a certain frequency. In order to realize a stable system, the capacitor must be limited to values such that the added zero does not interfere with the system response over the input frequency range. For this reason, small variations in the value of the ESR in an output capacitor can have a significant destabilizing effect on the entire system. A major goal of electronic system design is to avoid limiting circuit components to a precise value or range of values, allowing for easy replacement and substitution of components.\nAnother previous effort to stabilize control systems involves raising the quiescent current. The quiescent current, sometimes called the leakage current, is the portion of the input current that does not contribute to the load current. In other words, it is the current that the system consumes when no load current is being supplied. By raising the quiescent current, non-dominant poles in the system can be pushed to much higher frequency levels outside the system's normal operating range. A drawback of this stabilization method is that a higher quiescent current drains the batteries that power the system. For this reason many modern applications demand a low quiescent current for increased battery lifetime."}
{"text": "The present invention relates to memory cell layout of CMOS-type SRAM (Complementary Metal Oxide Semiconductor Static Random Access Memory) among semiconductor memory devices.\nThe SRAM memory cell that comprises six transistors and is made by a typical semiconductor CMOS process is widely used for system LSIs and so on.\nThe prior art layout pattern of the CMOS-type SRAM memory cell will be described below with reference to FIG. 8.\nThe prior art SRAM memory cell comprises nMOS drive transistors TN1 and TN2, nMOS access transistors TN3 and TN4, pMOS load transistors TP1 and TP2, polysilicon wires PL1, PL2, PL3 and PL4, wiring layers AL1 and AL2, and contacts CN1, CN2, CL1 and CL2.\nThe nMOS drive transistor TN1 and the nMOS access transistor TN3 are formed on an n-type diffusion region DN1 and the nMOS drive transistor TN2 and the nMOS access transistor TN4 are formed on an n-type diffusion region DN2. The pMOS load transistor TP1 is formed on a p-type diffusion region DP1 and the pMOS load transistor TP2 is formed on a p-type diffusion region DP2.\nGates of the nMOS drive transistor TN1 and the pMOS load transistor TP1 are connected to each other with the polysilicon wire PL1 and drains of them are connected to each other with the wiring layer AL1 via contact, thereby forming a first inverter (CMOS structure). Gates of the nMOS drive transistor TN2 and the pMOS load transistor TP2 are connected to each other with the polysilicon wire PL2 and drains of them are connected to each other with the wiring layer AL2 via contact, thereby forming a second inverter (CMOS structure). The wiring layer AL1 as an output node of the first inverter is connected to PL2 as an input node of the second inverter, and the wiring layer AL2 as an output node of the second inverter is connected to PL1 as an input node of the first inverter. Thereby a latch circuit for holding data is formed.\nA drain of the nMOS access transistor TN3 is connected to the wiring layer AL1 as the output node of the first inverter and source thereof is connected to a bit line (not shown) extending longitudinally via the contact CN1. A drain of the nMOS access transistor TN4 is connected to the wiring layer AL2 as the output node of the second inverter and source thereof is connected to another bit line (not shown) extending longitudinally via the contact CN2. Gates of TN3 and TN4 are connected to a word line (not shown) extending transversally via the contacts CL1 and CL2, respectively.\nWith such memory cell layout, long lateral distance allows a wide interval between two bit lines so that coupling capacitance between bit lines, which may cause a problem in micro process, can be reduced. Therefore, such memory cell layout is advantageous to speeding-up.\nNext, relationship between capability ratio of drive transistors and access transistors and stability of data holding in the SRAM memory cell will be explained with reference to FIGS. 9, 10 and 11.\nFIG. 9 shows a memory cell circuit diagram for evaluating stability of data holding. This circuit assumes the situation that the access transistors TN3 and TN4 turns on when the word line is in VDD level for a reading operation, and the bit line is raised to precharge level.\nFIG. 10 shows input/output characteristics of two inverter circuits (INV1, INV2) in the latch circuit.\nAin-Aout and Bin-Bout represent characteristics of INV1 and INV2, respectively and it is plotted so as to be Ain=Bout and Bin=Aout. Cross points P1 and P2 in this drawing are stable points and each point corresponds to memory data 0 or 1. In the plot, as area surrounded by two curved lines becomes larger, stability of data holding at P1 and P2 improves. Here, when driving capability of the access transistors TN3 and TN4 becomes greater than that of nMOS transistors (drive transistors) TN1 and TN2 in the inverter circuits, input/output characteristics of the inverter circuits change as shown in FIG. 11. The reason is that the access transistors transmit VDD level of the bit line to the latch nodes more easily, so that area surrounded by two curved lines becomes smaller. When noise voltage is applied into the memory cell having such characteristics, cross points are reduced to be only P2xe2x80x2 and therefore the memory cell can hold only either data. That is, in the case where data other than P2xe2x80x2 (i.e. P1xe2x80x2) is held, the data is destroyed. Thus, maintaining a constant ratio of access transistors to drive transistors in driving capability is important for holding memory cell data stably. Generally, driving capability of access transistors is set to be 50 to 70% of that of drive transistors.\nIn the prior art SRAM memory cell, channel width of drive transistors is set to be larger than that of access transistors, thereby generating a difference between them in driving capability.\nIn the prior art SRAM memory cell in which channel width of the drive transistors is set to be larger than that of the access transistors, thereby generating a difference between them in driving capability, the diffusion regions necessarily include some bent parts and end parts. For example, in FIG. 8, the bent parts that produce round-offs as shown by dashed lines DL3 and DL4 are generated by difference between the nMOS drive transistors TN1 and TN2 and the corresponding nMOS access transistors TN3 and TN4 in channel width.\nWith such layout, at the bent parts of the diffusion regions, finish pattern is rounded off as shown by dashed lines DL1, DL2, DL3 and DL4 in the figure. As a result, a problem arises that transistor width of the nMOS transistors TN1, TN2, TN3 and TN4 becomes larger than required. Moreover, at the end parts of the diffusion regions, finish pattern is retreated as shown by dashed-lines DL5 and DL6. As a result, there arises a problem of reduction in overlap margin of the p-type diffusion region with respect to the contact as well as variation in channel width of the pMOS transistors TP1 and TP2.\nFurthermore, system mounted on a semiconductor chip has increasingly become large scale. In connection with this, there is a tendency that the block of SRAM with a large scale in bit capacity is mounted on the chip. In order to meet these requests on the system side, it is desired to further reduce the size of SRAM memory cell. Although it is effective to use a MOS transistor with smaller channel width for the purpose of reducing cell size, such small-sized pattern is prone to undergo great variations in characteristics due to processing fluctuations. Therefore, reduction in cell size makes stable design by sufficient operational margin difficult. On the other hand, with recent micro process, it is more difficult to obtain desirable processed form and round-off or retreat of pattern tend to take place. Moreover, there often causes the phenomenon that even the same pattern form changes in finished form due to peripheral pattern form.\nTo suppress such changes of processed form, it has already been implemented to correct mask pattern in consideration of bend up and bend down of layout pattern concerned and in consideration of the peripheral layout pattern in recent micro process. Such process, however, is sensitive to apparatus used in semiconductor diffusion process and processing conditions. Further, correction value must be modified each time processing conditions in diffusion process are changed, adding a burdensome operation.\nThe present invention is made to solve the above-mentioned problems. With the memory cell layout of the semiconductor memory device according to the present invention, it is possible to lay out diffusion regions in linear shapes without any bent part by generating a difference between access transistors and drive transistors in driving capability without changing their channel width. As a result, processed form of diffusion areas of the SRAM memory cell is hard to change, thereby suppressing variations in characteristics of transistors so that transistors with narrow channel width can be used. An object of the present invention is to provide a highly integrated semiconductor memory device by use of such transistors with narrow channel width\nA semiconductor memory device according to one aspect of the present invention has a SRAM memory cell comprising: a first inverter including a first nMOS transistor and a first pMOS transistor; a second inverter including a second nMOS transistor and a second pMOS transistor; a third nMOS transistor; and a fourth nMOS transistor, wherein an input node of the first inverter is connected to an output node of the second inverter; and an input node of the second inverter is connected to an output node of the first inverter; either of drain and source of the third nMOS transistor is connected to the output node of the first inverter; the other of drain and source thereof is connected to a first bit line; and gate thereof is connected to a word line; either of drain and source of the fourth nMOS transistor is connected to the output node of the second inverter; the other of drain and source thereof is connected to a second bit line; and gate thereof is connected to the word line; a first diffusion region forming the first nMOS transistor and the third nMOS transistor and a second diffusion region forming the second nMOS transistor and the fourth nMOS transistor, respectively, are arranged in linear shapes without having any bent part, and current driving capability of the first and second nMOS transistors is higher than that of the third and fourth nMOS transistors.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that gate length of the third and fourth nMOS transistors is longer than that of the first and second nMOS transistors.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that threshold voltage characteristics of the third and fourth nMOS transistors are higher than those of the first and second nMOS transistors.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that a gate oxide film of the third and fourth nMOS transistors are thicker than that of the first and second nMOS transistors.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that driving voltage of the word line is lower than power supply voltage supplied to the first and second inverters.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that a third diffusion region forming the first pMOS transistor and a fourth diffusion region forming the second pMOS transistor are arranged in linear shapes without having any bent part and located in parallel to the first and second diffusion regions.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device further comprising: a third pMOS transistor formed adjacent to drain of the first pMOS transistor on the third diffusion region; and a fourth pMOS transistor formed adjacent to drain of the second pMOS transistor on the fourth diffusion region, wherein gates of the first nMOS transistor, the first pMOS transistor and the fourth pMOS transistor are connected in succession via a first polysillicon wire; and gates of the second nMOS transistor, the second pMOS transistor and the third pMOS transistor are connected in succession via a second polysillicon wire; and absolute value of threshold voltage of the third and fourth pMOS transistors is higher than power supply voltage supplied to the first and second inverters.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that n-type diffusion regions are provided on the opposite side of the diffusion region of the third pMOS transistor to the first pMOS transistor and on the opposite side of the diffusion region of the fourth pMOS transistor to the second pMOS transistor to fix potential of N-well region forming pMOS transistors.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that, in a memory cell array in which the SRAM memory cells are arranged in a grid pattern, a substrate contact region for fixing well potential is provided at regular intervals in the extending direction of the first and second diffusion regions; a diffusion region forming the substrate contact region is arranged in a linear shape without any bent part on the extension of the first diffusion region and the second diffusion region of the memory cell; and a fifth nMOS transistor is provided at the boundary of the memory cell and the substrate contact region wherein gate potential thereof is fixed so as not to become ON state.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that, in a memory cell array in which the SRAM memory cells are arranged in a grid pattern, a substrate contact region for fixing well potential is provided at regular intervals in the extending direction of the third and fourth diffusion regions;, a diffusion region forming the substrate contact region is arranged in a linear shape without any bent part on the extension of the third diffusion region and the fourth diffusion region of the memory cell; and a fifth pMOS transistor is provided at the boundary of the memory cell and the substrate contact region wherein gate potential thereof is fixed so as not to become ON state.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that the first pMOS transistor is formed on the third diffusion region; the second pMOS transistor is formed on the fourth diffusion region; and the first, second, third and fourth diffusion regions are arranged at equally spaced intervals in the longitudinal direction of connect line of transistor gate.\nA semiconductor memory device according to another aspect of the present invention is an above-mentioned semiconductor memory device characterized in that the first pMOS transistor is formed on the third diffusion region; the second pMOS transistor is formed on the fourth diffusion region; and the first, second, third and fourth diffusion regions are arranged so as to have same width measured in the longitudinal direction of connect line of transistor gate.\nAs a means of generating a difference in transistors\"\" capability without changing channel width, channel length of the access transistors is set to be longer than that of the drive transistors. This enables memory cell data to be held stably. In the prior art difference in capability has been generated from difference in channel width. Therefore, even if minimum width of diffusion region, which is possible in processing, can be applied to the access transistors, the drive transistors cannot be laid out at minimum width. According to the present invention, however, minimum width of diffusion region can also be applied to the drive transistors. In typical transistors, channel width is longer than channel length, and therefore reduction in channel width in the present invention probably can reduce cell size.\nHowever, in the prior art memory cell layout as shown in FIG. 8, because of its long sideways structure, longer channel length, which increases cell height, weakens the effect of reducing area by reduction in channel width of the drive transistor. Accordingly, in the present invention, as another means of generating a difference in transistors\"\" capability, threshold voltage of the access transistor is set to be greater than that of the drive transistor. Further, as other means, gate oxide film of the access transistor is set to be thicker than that of the drive transistor. By use of these means, the access transistor and the drive transistor have different capability while having the same channel length and width so that the effect of reducing area can be obtained. In the case that the above-mentioned means for changing characteristics of the access transistor are employed, it is necessary to add a supplementary process, suffering a slight increase in process costs. Therefore, as another means according to the present invention, power supply voltage supplied to a word line driving circuit for driving the gate of the access transistor is set to be lower than that supplied to a latch circuit in the memory cell. As a result, it is possible to generate a capability difference between the access transistor and the drive transistor having the same size and characteristics, thereby holding process costs down.\nAs for the p-type diffusion region, removing end parts that exist in the prior art layout is employed. The p-type diffusion regions are arranged in linear shapes, and for the separation of pMOS load transistor devices between vertically adjoining cells, a pMOS separating transistor having an absolute value of threshold voltage greater than power supply voltage and being in OFF state at any time irrespective of gate potential is formed. This suppress changes of processed form in channel width of the load transistor due to retreat of pattern that occurs at the end part in the prior art, and also allows a sufficient overlap of the diffusion region with respect to the contact in the drain of load transistor. Looking from the aspect of the n-type diffusion region, the p-type diffusion region always lies next to the n-type diffusion region as its peripheral layout form. In the prior art layout, as the p-type diffusion region does not exist partly next to the n-type diffusion region, the processed form of the n-type diffusion region changes a part and the change must be corrected on the mask data. According to the present invention, the processed form of the n-type diffusion region can be improved without any correction.\nFurthermore, fixing potential of N well forming the pMOS transistors by applying n-type diffusion on the opposite side of a pMOS separating transistor to load transistor is employed. Therefore, vacant region in the memory cell can be used effectively as substrate contact region. The substrate contact region for fixing potential of P well may be separately formed outside of the memory cell. Nevertheless, as potential of N well can be definitely fixed, latch-up resistance is improved. Moreover, a potential-fixed region of P well that formed separately can be set at larger intervals, thereby enabling area of the memory cell array to be suppressed.\nConsidering the layout of the whole memory cell array, it is necessary to provide the potential-fixed region of P well at regular intervals in the vertical direction of aligning memory cells. With the prior art layout, diffusion regions of memory cells above and under the potential-fixed region are necessarily separated. According to the present invention, by providing transistors whose gates are fixed in OFF state above and under the potential-fixed region both in n-type diffusion region and p-type diffusion region, the above-mentioned potential-fixed region of N well can be arranged in a linear shape in contiguous to the diffusion region of the memory cell. In this way, over the whole memory cell array, diffusion regions can be laid out in linear shapes without any bent part and partial separation, so that change in characteristics of transistors above and under the potential-fixed region can be suppressed.\nFurthermore, in the present invention, placing each diffusion region at the same interval is employed. As mentioned above, in recent micro process, finished processed form changes due to peripheral layout form. In the case of plural placement intervals, it becomes difficult to shape plural diffusion regions into desired interval and width. Maintaining a constant placement interval facilitates keeping forms of these diffusion regions constant. Similarly, placing each diffusion region at the same width is employed. It also facilitates keeping width and forms of these diffusion regions constant.\nThe novel features of the invention are set forth with particularity in the appended claims. The invention as to both structure and content, and other objects and features thereof will best be understood from the detailed description when considered in connection with the accompanying drawings."}
{"text": "Folding frames having four or more pivot points may be relatively easy to design. Folding frames having three pivot points exactly may be rather difficult to design but may result in functional structures that have the added benefit of being aesthetically pleasing to many eyes."}
{"text": "Surfboards may generally be described as water planing devices used to ride waves. The term “surfboard” may include boogie boards, wind surfing boards, and other hulled craft which are maneuvered by shitting the weight borne by the craft relative to the craft's center of gravity. Surfboards may be constructed of various lengths, width, shapes and thicknesses. Initially, surfboards had a single vertical fin located along the centerline of the surfboard at the rear that provided directional stability. Later designs added additional smaller fins of various sizes and shapes along the sides of the surfboard to improve the stability and maneuverability of the surfboards.\nIn terms of its design, a surfboard fin is analogous to a “fixed wing.” The surface curvature of conventional fins reflects architecture similar to that of a “fixed wing” aircraft. When the physics are applied to an aqueous environment, the “fixed wing” is termed a “hydrofoil” or sometimes simply a “foil.”\nHolding all other variables constant and for a constant hydrofoil speed, the velocity of water passing over the top surface of the hydrofoil wing is greater than that which passes below the flat bottom surface. With an increase in water velocity over the top of the hydrofoil wing, the pressure of the water over this surface is reduced when compared to the pressure below the hydrofoil wing. This difference in pressure tends to push the hydrofoil wing toward the side of lowest pressure of the pressure gradient. The velocity of the water over the top surface of the wing is forced to rise with increased velocity because (a) water is essentially incompressible and (b) the distance a water molecule must travel for a given linear distance is longer due to the curvature of the wing over the hydrofoil surface as compared to the straight line distance enjoyed by water molecules passing below the hydrofoil wing. The upward force generated through relative motion between the wing and the environment through which it travels, in this case water, is proportional to the velocity of the water over the wing surfaces and the surface area of the wing's bottom surface. The velocity of a surfboard is largely a function of wave height, wave velocity, and gravity.\nA reduction in lift may be expected from “turbulence” or “cavitation” that could theoretically develop over the wing surface. For example, turbulence may be created by high surface friction coefficients and disrupt desired laminar flow. Cavitation is a phenomenon whereby gasses actually come out of solution in the surrounding water mass due to significant pressure reduction. Similar to opening a carbonated beverage, and under the appropriate temperature and pressure conditions, air bubbles may form in the low pressure zones over the wing's surface. Cavitation may disrupt laminar flow and will also produce aquatic sounds as the bubbles collapse and return to solution.\nThe single vertical fin is generally speaking, a type of hydrofoil, whose function is to provide lift or some other force to the surfboard in reaction to its motion through the water. The single vertical fin is usually a symmetrical foil (a “50/50” foil) with both sides convex, which provides for even water flow on both sides of the fin such that a single vertical fin promotes stability and control. If the hydrofoil has a convex, top surface and a flat bottom surface, the velocity of water flow over the top surface of the hydrofoil is increased, thus creating a water pressure differential between the bottom and top surfaces and producing Sift or thrust in the direction of the pressure differential toward the area of lowest pressure.\nThe performance of a surfboard is affected by the design, placement, and number of fins affixed to the surfboard. For example, a fin may be defined by its dimensions: its base, its depth, its sweep, its flex, and its cant, and changes in these dimensions affect the performance characteristics of the surfboard. As for the hydrofoil effects, the fin may be the aforementioned symmetrical foil or a flat foil having one flat side and one convex side, which promotes maneuverability and fast transitions between turns. Other fins may combine the characteristics of a symmetrical foil and a flat foil in various proportions dependent on the desired performance characteristics of the surfboard. For optimum foil performance, flexibility in design characteristics is necessary, as well as the ability to modify these design characteristics as surfing conditions change.\nTherefore, as competitive surfers attempt more challenging maneuvers and ride bigger waves, there is need for improved fin designs that improve the lift, maneuverability, and other performance characteristics of the surfboards."}
{"text": "There are over nine million dairy cows in the United States and Canada, and over twenty million worldwide. The dairy industry is a very competitive marketplace, and the pregnancy status of the herd is critical to maximizing profits. It is estimated that a non-pregnant cow costs the industry approximately five dollars per day. An accurate, rapid test for determining the pregnancy status of a herd would have a very important economic impact on ranch or farm operations and would increase milk production of the dairies, resulting in increased profitability for the dairies.\nA number of antigens are known to be present in cows and sheep during pregnancy, and pregnancy has been evaluated by a variety of methods. Bovine Antigen Glycoprotein (U.S. Pat. No. 4,755,460, issued Jul. 5, 1988, and U.S. Pat. No. 4,895,804, issued Jan. 23, 1990) can be measured about 12-15 days after breeding. Early Pregnancy Factor (EPF) (U.S. Pat. No. 4,877,742, issued Oct. 31, 1989, and WO 00/51520, published Sep. 8, 2000) levels can be measured at about 20-40 days after breeding, such as with KEMS BioTest Ltd. (Littleton, Colo.) Animal Rapid Test for Bovine Pregnancy.\nInterferon-tau is produced by bovine trophoblast tissue between days 15-24 of bovine gestation and prevents luteolysis by suppressing endometrial PGF2α secretion. Interferon-tau induces or upregulates expression of a number of proteins in pregnant animals.\nProteins that are induced by IFN-T include granulocyte chemotactic protein (GCP-2) (WO 94/12537, published Jun. 9 1994 and Staggs, K. L. et al. [1998] “Complex Induction of Bovine Uterine Proteins by Interferon Tau” Biol. Reprod. 59:293-297), 2′,5′-oligoadenylate synthetase (Short, E. C. et al. [2001] “Expression of antiviral activity and induction of 2′,5′-oligoadenylate synthetase by conceptus secretory proteins enriched in bovine trophoblast protein-1” Biol. Repro. 44:261-268), β2-microglobulin (Vallet, J. L. et al. [1991] “A low molecular weight endometrial secretory protein which is increased by ovine trophoblast protein-1 is a β2-microglobulin-like protein,” J. Endocrinology 130:R1-R4), IFN regulatory factors 1 (IRF-1) and 2 (IRF-2) (Spencer, et al. [1998] Biol. Reprod. 58:1154-1162; and Binelli M. et al. [2001] Biol. Reprod. 64(2):654-665), GCP-2 (Teixeira, M. G. et al. [1997] Endocrine 6:31-37); and 1-8U, 1-8D, and Leu-13/9-27 (Pru, J. K. et al. [2001] “Pregnancy and Interferon -T Upregulate Gene Expression of Members of the I-8 Family in the Bovine Uterus” Biol. Reprod. 65:1471-1480; and Pru, J. K. [2000] “Regulation of bovine uterine proteins and prostaglandin F2a release by interferon-tau” Ph.D. Thesis, University of Wyoming). Leu-13 is the name of the protein encoded by the 9-27 gene. Cyclooxygenase-2 (COX-2) (Xiao, CW et al. [1998] “Regulation of COX-2 and prostaglandin F2a synthase gene expression by steroid hormones and IFN-T in bovine endometrial cells,” Endocrinol. 139:2293-2299 and Thatcher, W. W. et al. [2001] A Uterine-conceptus Interactions and Reproductive Failure in Cattle” Theriogenology 56:1435-1450) and PLA2 (Binelli, M. et al. [2000] “Interferon-tau modulates phorbol ester-induced production of prostaglandin and expression of cyclooxygenase-2 and phospholipase-A2 from bovine endometrial cells” Biol. Repro. 63:417-424) are also regulated by IFN-T.\nTeixeira, M. G. et al. (1997) “Bovine Granulocyte Chemotactic Protein-2 is Secreted by the Endometrium in Response to Interferon-tau,” Endocrine 6(1):31-37 report that bovine 1-8 transcripts were detected on Days 15 and 18 of pregnancy and were absent on Day 12 of pregnancy and during the estrus cycle. Bovine 1-8 gene family members are not known to be secreted. This reference also reported that polyclonal antibodies to a GCP-2 peptide, were generated in sheep, and used to demonstrate that GCP-2 is secreted by cultured endometrial cells, representing Day 14 of the estrus cycle, when dosed with IFN-T. \nMx encodes a monomeric GTPase and is induced by IFN-T (Ott, T. L. et al. [1988] “Effects of the Estrous Cycle and Early Pregnancy on Uterine Expression of Mx Protein in Sheep (Ovis aries)” Biol. Reprod. 59:784-794). In Ott et al. (1998), ovine Mx protein was detected using a monoclonal antibody directed against the amino terminus of human MxA (1319.35.126, supplied by M. Horisberger, Novartis, Basel Switzerland) and a Super ABC Mouse/Rat Kit (Biomeda, Foster City Calif.). U.S. patent application Ser. Nos. 60/299,553 and 10/166,929 describe a method of determining pregnancy status of an animal by assaying the level of Mx and comparing it to the level of Mx in a non-pregnant animal. Mx protein was detected with ovine Mx peptide antiserum (#90618-2). Yankey, S. J. et al. (2001) “Expression of the antiviral protein Mx in peripheral blood mononuclear cells of pregnant and bred, non-pregnant ewes” J. of Endocrinology 170:R7-R11, describes the presence of Mx in peripheral blood mononuclear cells of pregnant ewes at Day 15 of pregnancy. Mx protein can also be used to detect viral infection (EP 0 725 081, published Aug. 7, 1996) using monoclonal antibodies to human Mx. Antibodies to human Mx and immunoassays for Mx have been described (Staeheli, P. and Haller, O. [1985] “Interferon-induced human protein with homology to protein Mx of Influenza virus-resistant mice” Mol. Cell. Biol. 5(8):2150-2153; Towbin H. et al. [1992] “A Whole Blood Immunoassay for the Interferon-Inducible Human Mx Protein” J. Interferon Res. 12(2):67-74; U.S. Pat. Nos. 5,869,264, issued Feb. 9, 1999; U.S. Pat. No. 5,739,290, issued Apr. 14, 1998; and U.S. Pat. No. 6,180,102 issued Jan. 30, 2001). Antibodies to mouse Mx are described in Staeheli, P. et al. (1985) Mol. Cell Biol. 5:2150-2153; Staeheli, P. et al. (1985) J. Biol. Chem. 260(3):1821-1825; and Horisberger, M. A. et al. (1985) J. Biol. Chem. 260(3):1730-1733. One of the monoclonal antibodies in Towbin (1992) is reported to react with other species' Mx proteins (mouse, rat, bovine, and porcine), in addition to human Mx.\nAnother IFN-T-induced protein is ubiquitin cross-reactive protein (UCRP), which was first identified in humans (Farrell, P. J. et al. [1979] Nature 279:523-525) and later characterized (Koran, B. D. [1984] “Interferon-induced Proteins” J. Biol. Chem. 259(23):14835-14839; Blomstrom, D. C. et al. [1986] J. Biol. Chem. 261:8811-8816; and Knight E. Jr. et al. [1988] J. Biol. Chem. 263:4520-4522). Human UCRP (hUCRP) and mouse UCRP encode proteins that are processed to 17 kDa but that migrate as 15 kDa on PAGE gels (Potter, J. L. et al. [1999] “Precursor processing of pro-ISG15/UCRP, an interferon-beta-induced ubiquitin-like protein” J Biol. Chem. 274:25061-25068). These proteins are similar to ubiquitin, and are upregulated by interferon (IFN), hence they are also known as interferon-stimulated gene 15 (ISG15). ISG15 is involved in the viral response and in the recognition of pregnancy (Bebington, C. et al. [1999] “Localization of Ubiquitin and Ubiquitin Cross-Reactive Protein in Human and Baboon Endometrium and Decidua During the Menstrual Cycle and Early Pregnancy” Biol. Reprod. 60:920-928, and Bebington, C. et al. [1999] “Ubiquitin Cross-Reactive Protein Gene Expression is Increased in Decidualized Endometrial Stromal Cells at the Initiation of Pregnancy” Molecular Human Reproduction 5(10):966-972). Like ubiquitin, ISG15 becomes covalently attached to targeted intracellular proteins via a C-terminal LRGG amino acid sequence. Proteins that are coupled to ubiquitin often are degraded through the 26 S proteasome (Baboshina, O. V. [1996] “Novel multiubiquitin chain linkages catalyzed by the conjugating enzymes ESEPF and RAD6 are recognized by 26 S proteasome subunit 5,” J Biol. Chem. 271:2823-2831). Ubiquitin is conjugated to other proteins by E2-conjugating enzymes (Tanaka, K. et al. [1998] “The ligation systems for ubiquitin and ubiquitin-like proteins” Mol. Cell 8:503-512).\nThe 17 kDa bovine analog of hUCRP (ISG15) was identified as bovine UCRP (bUCRP) or ISG17 (Austin, K. J. et al. [1996] “Ubiquitin Cross-Reactive Protein is Released by the Bovine Uterus in Response to Interferon During Early Pregnancy,” Biol. Reprod. 54:600-606; Austin, K. J et al. [1996] “Complementary Deoxyribonucleic Acid Sequence Encoding bovine Ubiquitin Cross-Reactive Protein,” Endocrine 5(2):191-197; and Perry, D. J. et al. [1999] “Cloning of Interferon-Stimulated Gene 17: The Promoter and Nuclear Proteins That Regulate Transcription,” Molecular Endocrinology 13:1197-1206). ISG17 becomes covalently linked to targeted intracellular proteins, is released from endometrial cells, and may function as a paracrine modulator. Unlike ISG15, ISG17-conjugated proteins continue to accumulate rather than be degraded. Two of the 1-8 gene family members, bovine 1-8U and bovine Leu-13, have high homology with the E2-conjugating enzymes, and they retain critical amino acids for function. It has been suggested that they may function by conjugating ISG17 to proteins.\nA normal bovine estrus cycle is about 21 days in length. ISG17 has been detected by Day 15 of pregnancy. It continues to increase to Day 17, and remains high through Day 26 (Hansen, T. R. et al. [1997] “Transient Ubiquitin Cross-Reactive Protein Gene Expression in the Bovine Endometrium,” Endocrinology 138(11):5079-5082; and Spencer, T. E. et al. [1999] “Differential Effects of Intrauterine and Subcutaneous Administration of Recombinant Ovine Interferon Tau on the Endometrium of Cyclic Ewes,” Biol. Reprod. 61:464-470). ISG17 was not detectable above background during the estrus cycle of non-pregnant bovine.\nOne ISG17 function is to become cross-linked to cellular proteins, as does ubiquitin. Conjugation of ISG17 to endometrial cytosolic proteins was observed by Western Blotting using a polyclonal antibody to an ISG17 polypeptide (Johnson, G. A. et al. [1998] “Pregnancy and Interferon-Tau Induce Conjugation of Bovine Ubiquitin Cross-Reactive Protein to Cytosolic Uterine Proteins,” Biol. Reprod. 58:898-904). The peptide used to generate the polyclonal antibodies corresponds to amino acids 82 to 99 of ISG17. This polypeptide was chosen because it had a high antigenic index, 78% identity with ISG15, and low identity (22%) with ubiquitin. Attempts to use the antiserum to develop a pregnancy test met with limited or no success (Pru, J. K. [2000] “Regulation of bovine uterine proteins and prostaglandin F2a release by interferon-tau” Ph.D. Thesis, University of Wyoming, Appendix 1, page 1). Another antibody which has been utilized in the study of ISG17 is monoclonal antibody 5E9 (Pru, J. K. [2000] “Regulation of bovine uterine proteins and prostaglandin F2a release by interferon-tau” Ph.D. Thesis, University of Wyoming, Appendix 1).\nThe Johnson polyclonal antibody to ISG17 amino acids 82-89 was also used to study ISG17 induction by IFN-T by Western blotting (Staggs, K. L. et al. [1998] “Complex Induction of Bovine Uterine Proteins by Interferon Tau,” Biol. Reprod. 59:293-297).\nISG17 also can induce expression of IFN-T in peripheral blood mononuclear cells (PMBCs) (Pru, J. K. et al. [2000] “Production, Purification, and Carboxy-Terminal Sequencing of Bioactive Recombinant Bovine Interferon-Stimulated Gene Product 17,” Biol. Reprod. 63:619-628).\nOvine UCRP (oUCRP) has been cloned (Charleston, B. and Stewart, H. J. [1993] “An interferon-induced Mx protein: cDNA sequence and high level expression in the endometrium of pregnant sheep,” Gene 137:327-331). Ovine UCRP is reported to be detectable by Day 13, and to remain high through Day 19 of ovine pregnancy (Johnson, G. A. et al. [1999] “Expression of the Interferon Tau Inducible Cross-Reactive Protein in the Ovine Uterus,” Biol. Reprod. 61:312-318). Western blotting of oUCRP was performed using a polyclonal antibody to human UCRP.\nOther factors, in addition to IFN-T, may be responsible for the induction of UCRP (Johnson, G. A. et al. [2000] “Interferon-tau and Progesterone Regulate Ubiquitin Cross-Reactive Protein Expression in the Ovine Uterus,” Biol. Reprod. 62:622-627.\nEstrone sulfate was found to be increased around day 50 in bovine peripheral blood. (Hirako, M. and Takahashi, H. [2000], “Oestrone sulfate commences an increase around 50 days of gestation in bovine peripheral blood,” Reprod. Fertil. Dev. 12(7-8):351-354. Estrone sulfate analysis in urine or serum after Day 100 has also been used to confirm pregnancy (Holdsworth et al. [1982] J. Endocrin. 95:7-12 and Warnick et al. [1995] Theriogenol. 44:811-825).\nPSP60 is disclosed in Mialon, M. M., et al. (1993), “Peripheral concentration of a 60-kDa pregnancy serum protein during gestation and after calving and in relationship to embryonic mortality in cattle,” Reprod. Nutr. Dev. 33(3):269-82, to be present in peripheral blood from day 27 after artificial insemination until and beyond the end of pregnancy. Mialon, M. M., et al. (1994), “Detection of pregnancy by radioimmunoassay of a pregnancy serum protein (PSP60) in cattle,” Reprod. Nutr. Dev. 34(1):65-72 discloses that testing 349 cows for PSP60 28, 35, 50 and 90 days post-insemination gave accurate results compared with other known tests. Patel, O. V., et al. (1998), “Effect of stage of gestation and fetal number on plasma concentration of a pregnancy serum protein (PSP-60) in cattle,” Res. Vet. Sci. 65(3):195-199 discloses that PSP60 increased from day 20 post-oestrus to 20 days pre-partum.\nPregnancy-associated glycoprotein 1 (PAG-1) is disclosed in Xie, S., et al. (1991), “Identification of the major pregnancy-specific antigens of cattle and sheep as inactive members of the aspartic proteinase family,” Proc. Nat'l Acad. Sci. USA 88(22):10247-10251. This article teaches that pregnancy in cattle and sheep can be diagnosed by the presence of this conceptus-derived antigen in maternal serum. Zoli, A. P., et al. (1992), “Radioimmunoassay of a bovine pregnancy-associated glycoprotein in serum: its application for pregnancy diagnosis,” Biol. Reprod. 46(1):83-92 discloses a double-antibody radioimmunoassay for bovine PAG-1 which was detected in maternal peripheral blood beginning at day 22 of pregnancy and increasing progressively to day 270, and becoming undetectable by day 100 postpartum. Xie, S. et al. (1997), “The diversity and evolutionary relationship of the pregnancy-associated glycoproteins, an aspartic proteinase subfamily consisting of many trophoblast-expressed genes,” Proc. Nat'l Acad. Sci. USA 94(24):12809-12816, teaches that cattle, sheep and probably all ruminant artiodactyla possess up to 100 or more pregnancy-associated glycoprotein genes, many of which are placentally expressed. Szenci, O. et al. (1998), “Evaluation of false ultrasonographic diagnoses in cows by measuring plasma levels of bovine pregnancy-associated glycoprotein 1,” Vet. Rec. 142(12):304-306 taught that this antigen showed that before day 31 ultrasonographic scanning was not very sensitive because six of the 30 calving cows were incorrectly diagnosed as non-pregnant. 0.5 ng/ml was used as the cut-off point to determine pregnancy. Pregnancy Associated Glycoproteins (PAGs) can also be detected during early pregnancy (WO 99/47934, published Sep. 23, 1999). Szenci, O. et al. (1998) “Comparison of Ultrasonography, Bovine Pregnancy-Specific Protein B, and Bovine Pregnancy-Associated Glycoprotein 1 Tests for Pregnancy Detection in Dairy Cows” Theriogenology 50:77-88, describes a comparison of bovine pregnancy tests for days 26 to 58 after artificial insemination (AI). Green, J. et al. (2000), “Pregnancy-associated bovine and ovine glycoproteins exhibit spatially and temporally distinct expression patterns during pregnancy,” Biol. Reprod. 62(6):1624-1631, discloses that pregnancy-associated glycoproteins in sheep and cows are expressed in the trophectoderm or binucleate cells. Those expressed predominantly in bovine binucleate cells are expressed weakly if at all by day 25 placenta, but are present at the middle and end of pregnancy. Others, such as PAG-4, -5 and -9 are present at Day 25 and at earlier stages. Roberts, R. M., et al. (1995), “Glycoproteins of the aspartyl proteinase gene family secreted by the developing placenta,” Adv. Exp. Med. Biol. 362:231-240, teaches that pregnancy in cattle and sheep can be diagnosed by the presence of placentally-derived antigens (pregnancy-associated glycoproteins or PAG-1) in maternal serum soon after implantation begins at about day 20 following conception.\nPregnancy-specific Protein B (PSPB) is disclosed in U.S. Pat. Nos. 4,554,256, issued Nov. 19, 1985; 4,705,748, issued Nov. 10, 1987; European Patent No. 0132750, published Feb. 13, 1985; and Sasser, R., et al. (1986), “Detection of pregnancy by radioimmunoassay of a novel pregnancy-specific protein in serum of cows and a profile of serum concentrations during gestation,” Biol. Reprod. 35(4):936-942. Serum concentrations of PSPB exceeded 1 ng/ml by 30 days post-breeding and increased gradually through three months, six months, and nine months of gestation, declining steadily to less than 78 ng/ml by 21 days postpartum. PSPB could be measured in most cows by 24 days after breeding. Szenci, O. et al. (1998), “Comparison of ultrasonography, bovine pregnancy-specific protein B, and bovine pregnancy-associated glycoprotein 1 tests for pregnancy detection in dairy cows,” Theriogenology 50(1):77-88, teaches that at days 26 to 58 after artificial insemination, pregnancy testing with PSPB-diagnosed pregnant cows as accurately as measuring of PAG-1 or ultrasound; however, there were fewer false positive diagnoses with the PSPB test than the PAG-1 test. PSPB has also been tested in llamas (Drew, M. I. et al. [1995] “Pregnancy determination by use of pregnancy-specific protein B radioimmunoassay in llamas” JAVMA 207(2):217-219); deer (Willard, S. T. et al. [1998] “Early pregnancy detection and the hormonal characterization of embryonic-fetal mortality in fallow deer” Theriogenology 49:861-869; and sheep (Willard, J. M. et al. [1995] “Detection of fetal twins in sheep using a radioimmunoassay for PSPB” J. Anim. Sci. 73:960-966) for detection of twins. PSPB is also detectable after calving (Kiracofe, G. H. et al. [1993] “PSPB in serum of postpartum beef cows” J. Anim. Sci. 71:2199-2205). Polyclonal antibodies against PSPB are described in U.S. Pat. No. 4,705,748 and Humblot et al. (1988), “Pregnancy-specific protein B, progesterone concentrations and embryonic mortality during early pregnancy in dairy cows,” Reprod. Fertil. 83(1):215-223.\nProgesterone is an antigen which is present throughout pregnancy. Progesterone levels have been measured in milk or blood samples collected from cattle after 22-24 days, such as offered at Rocky Mountain Instrumental Laboratories Inc. (Fort Collins, Colo.), but measurements of progesterone in milk at days 18-22 yield unacceptably high rates of false positives (Oltenacu et al. [1990] J. Dairy Sci. 73:2826-2831 and Markusfeld et al. [1990] Br. Vet. J. 146:504-508). Moriyoshi, M. et al. (1996), “Early pregnancy diagnosis in the sow by saliva progesterone measurement using a bovine milk progesterone qualitative test EIA kit,” J. Vet. Med. Sci. 58(8):737-741 discloses that pregnancy could be diagnosed 17-24 days after last mating in sows. Polyclonal antibodies to progesterone are commercially available from many different sources including Research Diagnostics, Inc., Flanders, N.J., and are described in Humblot, F., et al. (1988) “Pregnancy-specific protein B., progesterone concentrations and embryonic mortality during early pregnancy in dairy cows,” Reprod. Fertil. 83(1):215-223. Monoclonal antibodies to progesterone are available commercially through OEM Concepts, Tom's River, N.J.\nJohnson, G. A. et al. (1998) “Pregnancy and Interferon-Tau Induce Conjugation of Bovine Ubiquitin Cross-Reactive Protein to Cytosolic Uterine Proteins,” Biol. Reprod. 58:898-904, discloses polyclonal antibodies to ISG17. The peptide used to generate the polyclonal antibodies corresponds to amino acids 82 to 99 of ISG17, LVRNDKGRSSPYEVQLKQ. This polypeptide was chosen because it had a high antigenic index, 78% identity with ISG15, and low identity (22%) with ubiquitin. Attempts to use the antiserum to develop a pregnancy test met with limited or no success (Pru, J. K. [2000] “Regulation of bovine uterine proteins and prostaglandin F2a release by interferon-tau” Ph.D. Thesis, University of Wyoming, Appendix 1, page 1). Another antibody which has been utilized in the study of ISG17 is monoclonal antibody 5E9 (Pru, J. K. (2000) “Regulation of bovine uterine proteins and prostaglandin F2a release by interferon-tau” Ph.D. Thesis, University of Wyoming, Appendix 1). U.S. patent application Ser. No. 60/393,615 discloses cDNAs believed to be associated with early bovine pregnancy.\nPrior bovine pregnancy tests have tested only single antigens. However, false positives may occur when single antigens are tested, since positive test results may occur for these antigens when certain viruses are present. Some antigens such as progesterone are present in lactating cows. Thus a test is needed which will reliably determine bovine pregnancy with minimal false positive results.\nMethods of making assay devices are described in Millipore's Short Guide for Developing Immunochromatographic Test Strips (2nd ed). Other assay devices and methods are described in U.S. Pat. Nos. 4,313,734, 4,376,110, 4,435,504, 4,486,530, 4,703,017, 4,740,468, 4,855,240, 4,954,452, 5,028,535, 5,075,078, 5,137,808, 5,229,073, 5591645, 5,654,162, 5,798,273, and in EP 0810436A1, and WO 95/16207. Assay devices containing more than one test strip are described at the Unitec, Inc. website.\nIn cows, the estrus cycle is about 21 days. To determine when a cycling cow is ready for breeding, the cow can be observed for behavioral estrus. Alternatively, a cow can be induced or forced into estrus with effective hormone therapies. Estrus of an entire herd can be synchronized (U.S. Pat. Nos. 3,892,855 issued Jul. 1, 1975, and 4,610,687 issued Sep. 9, 1986). Estrus synchronization, or preferably ovulation synchronization, is used in timed AI (TAI) breeding programs. TAI breeding programs involve precise estrus synchronization which allows for timed breeding without monitoring for behavioral estrus. Examples of methods for forcing estrus include U.S. Pat. No. 5,589,457 (issued Dec. 31, 1996), Ovsynch (Pharmacia Animal Health, Peapack, N.J.), Cosynch, Select Synch, Modified Select Synch, MGA/PGF, and Syncro-Mate-B. Such methods typically employ hormones such as prostaglandins, e.g. PGF2α (Lutalyse7, Pharmacia Upjohn, Peapack, N.J.; Bovilene7, Syntex; Animal Health, Des Moines, Iowa; and Estrumate7 Haver Lockhart, Shawnee, Kans.), and gonadotropin-releasing hormone (GnRH). Ovsynch involves a GnRH injection followed by a prostaglandin injection one week later, followed by a second GnRH injection 48 hours later. Insemination is ideally then performed at 12-18 hours, preferably about 16 hours, after the second GnRH injection. Ovsynch is maximally effective when implemented between Days 18-20 of a 20-day bovine estrus cycle (Thatcher, W. W. et al. [2000] “New Strategies to Increase Pregnancy Rates” at the website nab-css.org/education/Thatcher.html. Presynch (Pharmacia Animal Health, Peapack, N.J.) can be used to synchronize heifers before implementing Ovsynch. Presynch involves two prostaglandin injections. Some of the above-mentioned methods are also used on non-cycling cows to induce cycling, such as in lactating dairy cows. After precise estrus synchronization, animals need not be monitored for behavioral estrus and may be bred by appointment. Some animals may need estrus presynchronization before estrus synchronization. Melengestrol acetate (MGATM) in feed (Imwalle, D. B. et al. (1998) “Effects of melengestrol acetate on onset of puberty, follicular growth, and patterns of luteinizing hormone secretion in beef heifers” Biol. Repro. 58:1432-1436) or implants (U.S. Patent Publication No. 2001/0041697, published Nov. 15, 2001) can be used for presynchronizing estrus in heifers. Resynch is a program whereby animals are synchronized and bred, and then those animals that are determined to be open (not pregnant) are again synchronized and rebred.\nProstaglandin alone has been administered sequentially or simultaneously with artificial insemination to reduce the number of insemination administrations per herd required for achieving pregnancy (WO 02/04006, published Jan. 17, 2002).\nProstaglandin can be used as a single injection. An injection of about 2-5cc of Lutalyse (prostaglandin PGF2α) will induce an animal with a mature corpus luteum to come into estrus in about 48-96 hours. Cattle typically have a functional corpus luteum during Days 5-18 of the cycle (Estrus Synchronization of Cattle, Publication F-3163, Oklahoma Cooperative Extension Service, Oklahoma State University). Animals induced into estrus can be bred at 2-5 days following a prostaglandin injection. Single injection prostaglandin programs are often used with estrus synchronization, corpus luteum palpation, or behavioral heat detection because only animals in certain stages of the estrus cycle will respond by going into estrus. Breeding by appointment with a standard prostaglandin program has not been recommended because the interval from injection to estrus varies depending on the stage of the cycle when prostaglandin is administered. For example, if a cow is at cycle Day 7-8 or Day 15-17, timed AI can be performed at about 72-80 hours after the injection (O'Connor, M. L. discussion found at the website inform.umd.edu/EdRes/topic/AgrEnv/ndd/reproduce and das.psu.edu/reproduction/check/pdf/synchron.pdf. A risk of using prostaglandin injection for forcing estrus is that prostaglandin can cause abortion when given to pregnant animals. Estrus and ovulation synchronization allows cattle managers to concentrate heat detection efforts in a relatively short period of time or allows for TAI, which requires no heat detection (see the website ianr.unl.edu/beef/g741.htm).\nThere is a need in the art to determine pregnancy status during the breeding of livestock. In cattle, conception rates are low (Streenan and Diskin, Eds. [1986] Embryonic Mortality in Farm Animals, Martinus Nijhoff Publishers, 1-11) and spontaneous abortion rates are high, making pregnancy/non-pregnancy determination and rebreeding/inseminating important management tools. Particularly there is a need to determine pregnancy/non-pregnancy status during the estrus cycle in which insemination occurs or the first estrus cycle after insemination so that animals that are not pregnant can be most economically rebred. This need is particularly strong when raising livestock such as cattle, especially on dairy farms.\nThere is a need in the art for tests that determine pregnancy, and particularly non-pregnancy, status of animals during the estrus cycle in which insemination occurs or during the first estrus cycle after insemination. Knowing which animals are non-pregnant allows efforts to be directed towards forcing non-pregnant animals into estrus and/or watching for signs of estrus, in preparation for insemination, to decrease the time an animal is not pregnant. Pregnancy is dependent, not only on conception/fertilization but also on maternal recognition of pregnancy during the critical period, which allows for implantation. Up to 40% of total embryonic losses are estimated to occur between Days 8 and 17 of pregnancy in cattle (Thatcher, W. W. et al. [1994] “Embryo Health and Mortality in Sheep and Cattle,” J. Anim. Sci. 72(Suppl. 3):16-30). In the absence of reliable pregnancy tests, the earliest time at which a non-pregnant animal can be identified is at the beginning of a new estrus cycle, by observation of behavioral estrus. Optimally, pregnancy/nonpregnancy status is determined towards the end of or after the critical period when pregnancy is maintained, Days 15-17 according to Binelli, M. et al. (2001) “Antiluteolytic Strategies to Improve Fertility in Cattle,” Theriogenology 56:1451-1463, but before the end of the first estrus cycle, Days 18-20, allowing timed artificial insemination programs to be maximally effective. This reference discloses that pregnancy/non-pregnancy status is optimally determined during Days 17-18.\nAdditional technology relating to pregnancy testing in cows and other animals is disclosed in U.S. Provisional Patent Application Nos. 60/377,987, 60/377,166, 60/380,043, 60/377,921, 60/377,165, 60/377,355, 60/377,829, and 60/380,042, all filed May 2, 2002.\nAll references cited herein are incorporated herein by reference in their entirety to the extent that they are not inconsistent with the disclosure herein. Citation of the above documents is not intended as an admission that any of the foregoing is pertinent prior art. All statements as to the date or representation as to the contents of these documents is based on subjective characterization of the information available to the applicant, and does not constitute any admission as to the accuracy of the dates or contents of these documents."}
{"text": "The invention relates to portable screen systems.\nPreviously, it has been known to provide portable screen systems for separating fine material from coarse material wherein the frame of the unit is lowered flush onto preferably flat level ground, and material to be processed which overflows beyond side ends of the screen builds up along end sides of the unit. A set of wheels are made movable relative to the frame from an operative position for transporting the apparatus to an inoperative position for resting the frame flush on the ground. Such a unit is shown, for example, in U.S. Pat. Nos. 4,197,194, Des. 263,836, 4,237,000, and 4,256,572.\nThere are significant disadvantages with such a system. If the ground onto which the frame is rested in flush fashion is not level and/or bumpy and undulating, the vibrating screen will not be level side-to-side, and the material placed on the vibrating screen will be inefficiently processed and tend to shift more to one side of the screen than the other and thus be unevenly distributed. Accordingly, with the aforementioned unit, the ground on which the frame is positioned in flush manner should be level and not bumpy or undulating.\nA further disadvantage is that during operation of the aforementioned unit, overflow material to the sides of the unit builds up around the wheels and the side walls and can interfere with removal of the unit from the site due to such build-up of overflow material. Also, build-up of material where the frame sits flush on the ground may impede removal during freezing weather if such material were to freeze and trap the frame.\nAnother disadvantage of the above described unit is that if the frame resting on the ground settles unevenly during operation, the unit may become tilted, particularly if a rock was in contact with one small portion of the frame."}
{"text": "1. Field of the Invention\nThe present invention relates to a printed circuit board (PCB) antenna and, more particularly, to a PCB antenna capable of receiving four operating bands.\n2. Description of Related Art\nRecently, wireless communication products have gradually become a part of regular living. For example, mobile communication devices such as cell phones, have advanced to the “Third Generation” (3G) while ‘bluetooth’ products, providing great flexibility of PC devices etc, are becoming commonplace. The design of modem wireless communication products places great emphasis compactness, versatility, portability and aesthetics. However, when integrating a wireless communication product with a contemporary antenna, the size and appearance of the antenna seriously detract from the aesthetics of the communication product. Furthermore, the antenna can receive only a single band signal, which is not satisfactory. Currently, commercial wireless applications are approaching maturity, especially in reference to the computer information industry such that change from a wired network to a wireless network is well under way. However, in response to 2.4 GHz˜2.5 GHz/5.15 GHz˜5.25 GHz/5.25 GHz˜5.35 GHz/5.725 GHz˜5.85 GHz frequency bands being opened by the global wireless local network market, it has become necessary to integrate a number of different band antennas into a single wireless communication product. This integration will occupy too much space in the communication products and reduce convenience and reliability in use.\nTherefore, it is desirable to provide an improved antenna to mitigate and/or obviate the aforementioned problems."}
{"text": "1. Field of the Invention\nThe present invention is generally directed to an improved apparatus and method for intermittent manufacturing nitrogen gas, separated from air by using a permeable membrane in an automatic and unattended process. Particularly, the improved method and device adds a gas flow measuring device in a gas conduit downstream from the permeable membrane. The gas flow measuring device indirectly monitors nitrogen gas purity by measuring normal gas flow and detects deviation from desired nitrogen purity levels by measuring changing gas flow from the permeable membrane. The nitrogen and vent gas flow rates from the permeable membrane relate to the nitrogen purity so monitoring either or both the nitrogen and/or vent gas flow indirectly monitors the nitrogen gas purity. The automated process allows non-skilled people to use nitrogen gas from an automatic nitrogen producing apparatus. These people need assurances that the nitrogen gas meets purity requirements.\n2. Description of the Prior Art\nNitrogen manufacture from air by separating the oxygen and nitrogen has been accomplished by selective absorbent materials, distillation of liquid air, and membrane separation. These processes produce nitrogen for industrial uses such as chemical manufacture, inert gas welding, purging of explosive environments prior to electric arc cutting or welding, and food preservation. Also, these processes are mostly continuous nitrogen production to the industrial process or to continually fill large storage containers. The prior art is referenced in U.S. Pat. No. 5,588,984.\nThis U.S. Pat. No. 5,588,984, describes a method and apparatus to intermittently manufacture and dispense nitrogen. Air is filtered, compressed, and enters a module containing a permeable membrane that selectively separates nitrogen from the air and discharges vent oxygen and other gases. Automated temperature and pressure monitors allow the permeable membrane to separate air components. A discharge hose permits use of the nitrogen product for a variety of intermittent applications including vending for inflation of tires, filling portable nitrogen vessels, and use in manufacturing processes.\nDuring the use of this apparatus, debris, moisture, or oil collecting in the membrane may affect the permeable membrane efficiency of separating nitrogen from air. The membrane may need cleaning or other possible malfunctions from the air compressor, temperature, or pressure controls would change the apparatus efficiency or nitrogen purity. In large production and continuous manufacturing of nitrogen for storage or commercial use, a direct purity analysis method uses oxygen analyzers to indicate and monitor nitrogen purity. These oxygen analyzers directly determine the percent of oxygen, using electrochemical fuel cells thereby relating nitrogen purity to oxygen purity. One oxygen analyzer, manufactured by Teledyne, uses an expensive replaceable cell, needs monthly calibration, and is warranted for only six months while having a one year life. The analyzer is not rugged and would present a reliability problem for intermittent nitrogen production. An information sheet concerning these analyzers is enclosed. Oxygen gas analyzers usually require monitoring by manufacturing personnel to maintain the purity of nitrogen in constant flow nitrogen manufacture.\nFor automatic, intermittent operation by users of nitrogen, unfamiliar with the apparatus, a method and device for signaling an alarm and/or automatic shutdown of the nitrogen manufacturing apparatus is needed if the nitrogen purity is unacceptable. To replace the direct oxygen purity analysis equipment, an indirect monitor device using simple, reliable, and long life flow indicators, pressure monitors, and/or flow switches are used in intermittent apparatus operation down stream of the permeable membrane to monitor the purity of nitrogen gas. These gas flow devices do not require maintenance/and last appreciably longer than one year."}
{"text": "The present invention relates to a data handling system having a redundant storage configuration. More particularly, the present invention relates to a data storage suitable for averaging workload and for matching contents between storage devices having duplicate data, respectively, and constituting a redundant storage arrangement.\nJapanese Published Patent Application No. 06-259336 discloses a technique for matching contents between a master extended storage device and a sub extended storage device both of which constitute a redundant storage configuration. The disclosed technique involves matching the order in which to perform store operations on the master and the sub extended storage devices in order to keep contents identical between the master and the sub extended storage devices. However, the disclosure contains no reference to fetch operations on the master and the sub extended storage devices."}
{"text": "1. Field of the Invention\nThe inventive arrangements relate generally to the field of projection television receivers and displays and more particularly to projection television receivers and displays that employ imagers such as liquid crystal on silicon imagers.\n2. Description of Related Art\nThere have been many new developments in various types of electronic displays and video imaging devices. One example of such technology is liquid crystal on silicon (LCOS). As is known in the art, an LCOS imager generally contains an array of row and column electrodes such that the pixels of the LCOS imager can be addressed by selection of these row and column electrodes.\nTypically, a video input signal is selectively fed to each of the column electrodes, and selection of a row electrode enables each cell corresponding with the pixels to be charged to a desired pixel voltage. This permits video to be written to each of the rows of pixels. The video input signal is transferred to the column electrodes from a bus and through a number of switches connected to the bus and the column electrodes. These switches remain closed only for brief periods of time. A particular cell remains lighted with the same intensity until the video input signal changes that cell thereby acting as a sample and hold. That is, the pixel does not decay, as is the case with the phosphors in a cathode ray tube. Notably, many imagers permit the row electrodes to be selected in a sequential fashion, and some permit the row electrodes to be selected in a non-sequential manner.\nCurrent LCOS imagers, however, suffer from a significant drawback known as column memory. As the video input signal is transferred to a column electrode and the switch through which the input signal is passing opens, a charge remains on the column electrode. Thus, when the next row electrode is activated, the charge that is left over from the previous charging of the column electrode remains on the column electrode until the switch is closed again to write video to the new row of pixels. This residual charge can result in scene content from the previously written row being displayed in the new row being written thereby causing a phenomenon known as “ghosting.” The ghosting effect can be particularly troublesome if rows are selected in a non-sequential manner, as the voltage levels on the column electrodes from the previous row selection may be significantly different from the current row selection. Thus, it is desirable to eliminate the ghosting effect without significantly increasing system costs or complexity."}
{"text": "In an automatic transaction machine such as an automated teller machine, generally configuration is made, as shown in FIG. 5A and FIG. 5B, such that a rear door 102 is unlocked and opened, and an opening-and-closing cover 101 provided for example to an operation panel is opened by removing screws from the rear so as to enable maintenance such as of internal units.\nA cover opening-and-closing support mechanism is provided at an automatic transaction machine such as an automated teller machine in order to facilitate opening of the heavy opening-and-closing cover 101 when performing maintenance. A gas spring 91 shown in FIG. 5B assists raising of the opening-and-closing cover 101. A gas spring stopper 92, described later, also locks the gas spring 91 such that the opening-and-closing cover 101 cannot be closed.\nIn order to close the opening-and-closing cover 101 after maintenance has been completed the cover opening-and-closing support mechanism is configured so as to close slowly when the gas spring stopper 92 is pressed and released.\nA cover opening-and-closing support mechanism of such a related automatic transaction device 100 is configured, as shown in FIG. 4, with a gas spring 91 and a gas spring stopper 92.\nThe gas spring 91 is configured with a cylinder section 91c and a rod section 91d. A gas spring pivot point 91a is provided to the gas spring 91 at the top end of the cylinder section 91c, this being on the side of the opening-and-closing cover 101. A gas spring pivot point 91b is provided at the bottom end of the rod section 91d, this being on the side which is fixed to the cabinet of the automatic transaction device 100.\nThe gas spring stopper 92 is configured from a stopper portion 92a, a limiter portion 92b, and a spring 93. The stopper portion 92a locks the gas spring 91. The limiter portion 92b is provided so that the stopper portion 92a does not contact the rod section 91d of the gas spring 91 and damage the rod section 91d. The spring 93 is a biasing member for biasing the gas spring stopper 92 towards the gas spring 91 as indicated by arrow A, so as to cause the stopper portion 92a to contact the gas spring 91 as shown in the intermittent lined Region a.\nSuch a gas spring is attached to the opening-and-closing cover in order for example to push against the weight of the opening-and-closing cover during opening and closing. However, the reaction force from the gas pressure decreases as the gas spring 91 ages. The gas spring stopper 92 is therefore provided to prevent the opening-and-closing cover 101 from sagging. A stopper profile is provided at the gas spring stopper 92 that utilizes the difference between the external profiles of the cylinder section 91c and the rod section 91d of the gas spring 91.\nExamples of biasing members include a tension coil spring, as shown in FIG. 4, a compression coil spring that similarly presses the gas spring stopper 92 in the arrow A direction, and torsion springs that are for example attached to the gas spring pivot point 91a and bias in the arrow A direction (see for example Japanese Patent Application Laid-Open (JP-A) No. 10-115340)."}
{"text": "A structure of a 3GPP LTE (3rd Generation Partnership Project Long Term Evolution; hereinafter, referred as “LTE”) system which is an example of a wireless communication system to which the present invention may be applied will be described.\nFIG. 1 illustrates a schematic structure a network structure of an evolved universal mobile telecommunication system (E-UMTS). An E-UMTS system is an evolved version of the UMTS system and basic standardization thereof is in progress under the 3rd Generation Partnership Project (3GPP). The E-UMTS is also referred to as a Long Term Evolution (LTE) system. For details of the technical specifications of the UMTS and E-UMTS, refer to Release 7 and Release 8 of “3rd Generation Partnership Project; Technical Specification Group Radio Access Network”.\nReferring to FIG. 1, the E-UMTS includes a User Equipment (UE), base stations (or eNBs or eNode Bs), and an Access Gateway (AG) which is located at an end of a network (E-UTRAN) and which is connected to an external network. Generally, an eNB can simultaneously transmit multiple data streams for a broadcast service, a multicast service and/or a unicast service.\nOne or more cells may exist for one BS. The cell provides a downlink or uplink transmission service to several UEs using any one of bandwidths of 1.25, 2.5, 5, 10, 15 and 20 MHz. Different cells may be set to provide different bandwidths. A BS controls data transmission or reception to or from a plurality of UEs. The BS transmits downlink scheduling information to a UE with respect to downlink (DL) data so as to inform the UE of time/frequency domain, coding, data size, Hybrid Automatic Repeat and reQuest (HARQ) associated information of data to be transmitted, or the like. The BS transmits uplink scheduling information to a UE with respect to uplink (UL) data so as to inform the UE of time/frequency domain, coding, data size, HARQ associated information used by the UE, or the like. An interface for transmitting user traffic or control traffic can be used between BSs. A Core Network (CN) may include the AG, a network node for user registration of the UE, or the like. The AG manages mobility of a UE on a Tracking Area (TA) basis. One TA includes a plurality of cells.\nWireless communication technology has been developed to reach the LTE based on Wideband Code Division Multiple Access (WCDMA), but demands and expectations of users and providers have continuously increased. In addition, since other aspects of wireless access technology continue to evolve, new advances are required to remain competitive in the future. There is a need for reduction in cost per bit, service availability increase, the use of a flexible frequency band, a simple structure and an open type interface, appropriate power consumption of a UE, etc."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus such as a copying machine, a facsimile machine, a printer and a multifunction machine, and, moreover, to a sheet stacking device stacking a sheet (recording medium) formed with an image, and a sheet processing device performing post processing of a sheet.\n2. Description of the Related Art\nA sheet processing device with the following configuration has been well known as a sheet processing device into which a sheet with an image formed in an image forming apparatus is conveyed. The sheet processing device has a buffer roller through which, when the sheet processing device receives sheets, which have been formed with an image, and have been discharged from an image forming apparatus main body, the received sheets are superimposed for temporary waiting before the sheets are conveyed to a post processing mechanism such as a stapling machine and a saddle stitching machine. That is, while a preceding sheet bundle is processed in a processing tray, first several sheets of the succeeding sheet bundle are on the buffer roller for temporary waiting. When the preceding sheet bundle, which has been processed, is discharged from the processing tray, the succeeding several sheets, which have been delivered from the buffer roller, are conveyed to the processing tray. A brief explanation of a sheet post-processing device with the above-described configuration will be given, referring to FIG. 9A through FIG. 9C.\nA plurality of sheets P1, P2, . . . are superimposed one on top of another and wrapped around a buffer roller 5 to form a wrapping path 32. For example, three sheets P1, P2, and P3 are superimposed one on top of another, delivered from the path 32 after temporary waiting, conveyed, and conveyed to a processing tray 101 through a discharge roller 7, bundle discharge rollers 180a, and 180b. When the rear ends of the sheets pass the discharge roller 7, the bundle discharge rollers 180a and 180b rotate in the reverse direction in such a way that the sheet bundle of three sheets P1, P2, and P3 is returned in the direction in which the sheets abut against a rear-end stopper 3 of the processing tray 101. Alignment is performed in such a way that the bundle discharge roller 180b is separated from the bundle discharge roller 180a just before the rear end of the sheet bundle abuts against the rear-end stopper 3 and the sheet bundle abuts against the rear-end stopper 3 by moving inertia. At this time, alignment in a direction perpendicular to the conveyance direction is performed, using aligning plates.\nWhen all the sheets of the first sheet bundle are aligned on the processing tray 101 in such a manner, a swinging guide 150 is lowered and the bundle discharge roller 180b sits atop the sheet bundle to perform stitching processing of the sheet bundle, and the like, using a processing machine such as a stapling machine indicated by a reference number 4 in FIG. 9A.\nAccording to the above-described procedures, a first plurality of sheets of the succeeding second sheet bundle are wrapped around the buffer roller 5 as a temporary accumulating unit for waiting until processing for the first sheet bundle is completed. Thereby, a high-speed image forming apparatus by which sheets are discharged from the main body of an image forming apparatus at a small interval may be realized. Incidentally, the varieties of the quality and the size of sheets have been further increased in recent years. But the sheet processing device shown in FIG. 9A through FIG. 9C may hardly treat sheets, for example, special paper such as coated paper, the surface of which is treated, thick one, and large-sized one.\nEven if these kinds of sheets may be surely aligned one by one, it is difficult to align a plurality of the sheets in a state in which the sheets are superimposed. That is, the rear ends of a plurality of the sheets with a special quality, or with a special sheet size is run into the rear-end stopper 3 on the processing tray 1. In this case, there is generated a state in which all the three sheets P1, P2, and P3 with a large coefficient of friction, such as that of coated paper, are not completely returned to the rear end of the stopper 3. Especially, it is serious that the sheet P2 such as the second sheet of the superimposed ones is incompletely or faultily returned, that is, the quality of the post processing is deteriorated, and, consequently, the productivity is reduced. Moreover, when a plurality of sheets such as a thick one, and a large-sized one are superimposed and moved, there is caused larger inertia than the one caused in a case in which one sheet is moved. Accordingly, there is a case in which non-aligning is caused, because the sheets are vigorously run into the rear-end stopper 3 and bound. Moreover, there is a possibility that the end portion of the sheet buckles, and is damaged."}
{"text": "Many computing systems today utilize multiple processing units, resulting in a computer architecture generally referred to as multiprocessing. Multiprocessing systems are often used for transaction processing, such as airline and banking systems. Transaction processing refers generally to a technique for organizing multi-user, high volume, on-line applications that provides control over user access and updates of databases. A transaction refers to the execution of a retrieval or an update program in a database management system. Transactions originating from different users may be aimed at the same database records. This situation, if not carefully monitored, may cause the database to become \"inconsistent\". Where all transactions are executed one after the other, the database will remain in a consistent state.\nTo maintain transaction indivisibility, either all database updates in a transaction processing system or none of the updates are applied to the database. In prior art systems, database management systems have retained a copy of the existing data whenever the transaction requests a database update. If the transaction proceeds to a stable point, all new updates are secured to the database by the database management system. If the transaction does not proceed to a stable point because a failure occurs before the updates are secured to the database, none of the updates are applied. The prior art systems typically use control messaging between various components in the transaction processing system to ensure that each database transaction request is processed once and only once. If this complex synchronization between these various components is not maintained, the data in the database can become inconsistent. However, the synchronization between these components is very complex, and causes undesirable system overhead inefficiencies.\nTherefore, it is desirable to provide a system for providing transaction indivisibility without the complex synchronization between various system components. The present invention provides a solution to this problem by providing a centralized system accessible by all components, where an indivisibility analysis can be performed for all source messaging, regardless of which terminal or associated host processor initiated the transaction. The present invention therefore provides a solution to the aforementioned and other problems, and offers other advantages over the prior art."}
{"text": "Conventionally, in the coating of an automobile body, midcoat coating, baking, overcoat base coating, overcoat clear coating and baking are sequentially performed on a coated product on which electrodeposition coating serving as an undercoat and baking are performed. In other words, since a total of three rounds of baking are needed, a reduction in the number of times baking is performed has been required in terms of reducing the amount of CO2 discharged and saving energy.\nHence, a coating method has been proposed in which, in the coating of an automobile body, baking is not performed after midcoat coating, wet-on-wet overcoat coating is performed, and a midcoat coating film and an overcoat coating film are baked and cured simultaneously (see, for example, Patent Document 1). This coating method is a so-called 3-coat 2-bake (hereinafter referred to as “3C2B”) coating method in which a total of three rounds of coating (undercoat, midcoat and overcoat) and a total of two rounds of baking are performed.    Patent Document 1: Japanese Unexamined Patent Application, Publication No. 2005-177631"}
{"text": "Enterprises employ a vast number of protocols for internal and/or external communications and information transfer in order to ensure reliability, security and compliance with particular policies. For example, internal and/or external communications and information transfer can include, but are not limited to instant messaging (IM), electronic mail (email), Internet Protocol (IP) telephony, web mail, web-browsing, text messaging over a network of two or more computers (or network connectable, processor-based devices), and the like. These electronic communication media are popular as they provide inexpensive, easy, point-to-point communication that is less intrusive than traditional techniques and/or disparate non-electronic communications. There is an abundance of other benefits, for example, instant messaging (IM) is an electronic communication that easily enables one-to-many communication, there is no need to synchronize participants and the content can be planned more easily, among other things. Unfortunately, these media have adversaries and/or protocols that threaten the convenience of and confidence in their use, namely spam, viruses, malware, etc.\nA variety of systems and techniques have been developed and employed to combat spam and malicious code related to electronic communication media. With an increase in malicious activity involving spam, viruses, malware, and the like, enterprises are continuously searching for efficient techniques to secure networks and respective data communications associated therewith. Thus, conventional techniques (e.g., content-based filters, IP address-based filters, etc.) are becoming ineffective in recognizing and blocking disguised spam messages in relation to electronic communications. The integrity of enterprises and data communications is imperative to success in which improvements are necessary to ensure security and/or protective techniques are less vulnerable, stronger, and more difficult to penetrate."}
{"text": "A person's “cognitive load” is the degree to which the person's working memory is engaged in processing information. The more working memory is used, the higher the cognitive load. The higher a person's cognitive load, the greater the chances that “distracted operating” will impact the person's performance in operating a piece of equipment.\nFor instance, a person driving an automobile while operating a hand-held electronic device (e.g., text messaging on a cellular device) will experience an elevated cognitive load as the person tries to operate both the automobile (the equipment) and the hand-held electronic device at the same time. This impact on performance can be generally referred to as one type of “distracted driving.”\nAccording to the National Highway Traffic Safety Administration (NHTSA), “distracted driving” is a dangerous epidemic on America's roadways, as evidenced by the fact that in 2009 alone, it is estimated that nearly 5,500 people were killed, and 450,000 people were injured in distracted driving crashes.\n“Distracted operating” isn't limited to distractions caused by the utilization of hand-held electronic devices while operating equipment. Many other activities, events and situations can elevate a person's cognitive load. For instance, other activities, events and/or situations which can elevate a person's cognitive load include, but are not limited to, stress, mood, grief, the person's physical and/or mental health, the person's age, the person's maturity, cigarette smoking, eating, drinking, emergencies, sleepiness, weather conditions, the presence of hazards, “multi-tasking” (e.g., operating the equipment while also utilizing social media, playing games, watching television, listening to music, and/or talking on the telephone).\n“Impaired operating” (e.g., under the influence of alcohol, under the influence of illegal drugs, under the influence of prescription drugs) is likewise a concern. For instance, the impact of alcohol on a person's performance in operating an automobile is well known.\nTherefore, a need exist for methods of, and apparatuses for, determining if a person operating equipment is experiencing an elevated cognitive load."}
{"text": "In vehicles using electric traction motors, alternating current (AC) motor drives are used to provide a requested torque to the motor shaft. In practice, the amount of torque produced by the motor is directly related (although not perfectly proportional) to the amount current provided to the motor. Therefore, by regulating and precisely controlling the input current to the electric motor, the amount of torque produced by the electric motor may be more accurately controlled. In response to a changing torque command, the motor torque response is smoother and/or faster when the amount of current provided to the electric motor is adjusted based on the torque command.\nFor purposes of efficiency, particularly in hybrid and/or electric vehicles, it is desirable to maximize the ratio of the output torque to the input motor current. However, in many systems, the input motor current is not directly controlled. For example, many electric motors are operated using pulse-width modulation (PWM) techniques in combination with an inverter (or another switched-mode power supply) to control the voltage across the motor windings, which in turn, produces the desired current in the motor. In response to a requested torque (or commanded torque), most prior art systems determine a desired input motor current for producing the requested amount of torque and utilize a closed loop control system to control the current through the motor windings and thereby regulate the amount of torque produced the motor. One or more sensors are used to measure the actual motor current, which is then compared to the desired input motor current. Based on the outcome of the comparison, the PWM commands for the inverter are adjusted to increase and/or decrease the voltage across the motor windings, such that the actual measured motor current tracks the desired input motor current.\nHowever, when a current sensor does not accurately measure the motor current, these closed-loop control systems can no longer effectively control the motor torque. For example, without accurate motor current information, the control system may cause the motor to produce insufficient torque, excessive torque, or varying or oscillating amounts of torque. Furthermore, as a preventative measure in some prior systems, in response to a current sensor error, the control system may cease providing current and/or voltage to the electric motor, or drastically reduce the amount of current and/or voltage provided to the electric motor. As a result, in prior art systems, when a current sensor error occurs, the use and enjoyment of a particular vehicle is adversely affected."}
{"text": "(a) Field of the Invention\nThe present invention relates to a variable valve apparatus.\n(b) Description of the Related Art\nA typical combustion chamber of an automotive engine is provided with an intake valve, for supplying an air/fuel mixture, and an exhaust valve, for expelling burned gas. The intake and exhaust valves are opened and closed by a valve lift apparatus connected to a crankshaft.\nA conventional valve lift apparatus has a fixed valve lift amount due to a fixed cam shape. Therefore, it is impossible to adjust the amount of gas that is introduced or exhausted. However, valve timing and amount of lift should ideally be optimized for different driving speeds.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to a sealed container, a manufacturing method therefor, a gas measuring method, and a gas measuring apparatus for implementing the gas measuring method. More specifically, the invention relates to a sealed container used for a flat panel display, a manufacturing method for the sealed container, a gas measuring method used for measuring a gas rate of an emission gas, a leakage gas, or the like or measuring a life of a getter, and a gas measuring apparatus for implementing the gas measuring method.\n2. Related Background Art\nExamples of self-light emitting flat panel displays include a plasma display, an EL display device, and an image display device using an electron beam. An image display device using a sealed container that maintains its inside to a lower pressure than the atmospheric pressure is represented by a cathode ray tube (hereinafter, referred to as “CRT”) of a television set, but devices and apparatuses including the plasma display and a flat panel display using an electron beam also utilize the sealed container that has a pair of plates and maintains its inside to a lower pressure than the atmospheric pressure. Currently, there are increasing demands for the display devices to have a larger screen and a higher definition, and there are ever-growing needs for the self-light emitting flat panel displays.\nSuch image display devices face a major problem of an image display life. This is because, while having a gas source that may be hit by electrons and ions, the image display device must maintain a high vacuum for as long as several tens of thousands of hours by limited exhaust means, making it necessary for electron radiation from an electron source to be conducted in a stable manner over a long period of time. The radioactivity of the electrons from the electron source is largely influenced by an emission gas inside the image display device. For example, the CRT may involve a problem of damage caused by Ar (JP 10-269930 A).\nAccordingly, it is necessary to grasp types of gases causing damage to an electron source in an operation state and a gas generation rate (gas emission from a member) to reduce the damage to the electron source.\nFurther, in order to maintain a pressure inside a panel by the limited exhaust means, it is necessary to exhaust the emission gas emitted from the member. As the exhaust means, a barium getter is conventionally known, and almost all of its basic properties have become apparent. However, a gas absorbing power of the barium getter inside an actual panel is hard to estimate from the basic properties. This is because the absorbing power of a getter film largely differs according to a fine structure of the getter film inside the panel, the amount and type of the emission gas inside the panel (generation of a reaction product), and the like. Therefore, the absorbing power of a getter inside an actual panel can be only directly measured with respect to a subject panel.\nAccordingly, as a method of measuring a life of an image display device, it is a problem of urgency to establish a method of measuring a life of a getter, in which an influence of a gas exerted to a device when an image is displayed is evaluated (an emission gas rate is accurately measured for each type of gas) while a vacuum state of the image display device is maintained.\nOn the other hand, known as a conventional gas measuring method is a method of measuring a gas partial pressure using a quadrupole mass spectrometer (Q-Mass) as a mass spectrometer for analyzing gases inside a vacuum apparatus and a process chamber (JP 2952894 B).\nProposed as a method of measuring an emission gas rate and an adsorption gas rate for each gas is a measuring method using a partial pressure gauge provided to each of two chambers that are connected to each other through an orifice (JP 05-072015 A). Also, for a CRT, plural methods of measuring an emission gas rate and an adsorption gas rate are proposed as the method of measuring a life of a getter. Examples of the proposed plural methods include: a method of heating a CRT to 150° C. to 250° C. and measuring an emission gas rate while cooling the CRT (JP 07-226159 A); a method of measuring a gas absorbing power of a getter film after the CRT is caused to run for a predetermined period of time, calculating an amount of an emission gas from a built-in member of the CRT, and estimating a long-term life of a getter based on the calculated amount (JP 10-208641 A); and a method of finding a relationship between an amount of a getter and a life of a CRT by setting the amount of the getter to a small amount (JP 2000-076999 A).\nFurther, JP 2000-340115 A discloses a manufacturing method for an image display device in which a manufacturing process is performed while a state of an atmosphere is being monitored by using an orifice having a known conductance and installed in part of an exhaust channel of a manufacturing apparatus for vacuum pumping.\nAccording to the gas measuring methods disclosed in JP 2952894 B and JP 05-072015 A, a gas measurement is performed by placing a measurement sample inside a vacuum chamber and using a mass spectrometer, enabling the measurement for each type of gas. Particularly in JP 05-072015 A, a vacuum chamber having an orifice is used, enabling the measurement of an emission gas rate for each type of gas as well. However, it is difficult to place a large apparatus such as a flat panel display inside the vacuum chamber for the measurement. If the measuring apparatus is manufactured to be adapted for such a large apparatus, a huge manufacturing cost is required, making it hard to implement such arrangement.\nThe gas measurement for a CRT has long been performed. However, in JP 07-226159 A, a mass spectrometer is not used for the gas measurement, thereby making it impossible to measure an emission gas rate for each type of gas, and a gas to be adsorbed to a getter cannot be supplied, thereby making it impossible to accurately evaluate a life of a CRT. Further, in JP 10-208641 A, there are included an orifice and a total pressure gauge for measuring an emission gas rate, and a gas supply system for measuring a gas adsorbing power of a getter. However, a mass spectrometer is not used for a partial pressure measurement, thereby making it impossible to measure an emission gas rate for each type of gas. Also, it is possible to supply to the CRT a gas to be adsorbed to a getter through the orifice at a constant rate. However, lack of a chamber for adjustment of a pressure makes it difficult to adjust a pressure of the supplied gas, resulting in a long-time measurement. Further, according to the method of JP 2000-076999 A, which serves to measure the relationship between an amount of a getter and a life of a CRT by setting the amount of the getter to a small amount, the measurement requires a long period of time, and the gas measurement cannot be performed for a type of gas that is actually generated in the CRT. Therefore, it is difficult to accurately predict the life of the CRT.\nThe manufacturing method for an image display device disclosed in JP 2000-340115 A is suitable for a gas measuring method during the manufacturing, but is difficult to use as a gas measuring method for an image display device that has become a vacuum container.\nAlternatively, as the gas measuring method for a CRT that has been manufactured, there is a method in which a hole is opened by a punch when a pipe for a measurement is connected to a funnel of the CRT.\nHowever, according to this method, in the case of an apparatus using a thin glass plate such as a flat panel display, a crack easily develops, increasing the possibility of generating a leak."}
{"text": "Handheld computers are common electronic devices primarily functioning to provide a user with mobile computing services. Early handheld computers were limited to basic calendaring and phonebook applications. Now handheld computers are capable of nearly the same tasks as desktop or laptop computers, although memory capacity and processing speed may be limited.\nHandheld computers are often used in conjunction with a docking station to synchronize data files with a host computer and recharge the batteries of the handheld computer. These docking stations facilitate the connection of the handheld and host computer. A cradle or wired plug provides the electrical and structural elements of the docking station to dock a handheld computer. Once docked, data connectivity is established between the host and handheld computer. Optionally, an electrical connection is established to recharge the handheld computer's batteries. Other handheld computers utilize infrared or radio, particularly cellular telephone signals, to perform data exchanges with a host computer."}
{"text": "This invention relates to a reference clock architecture for an integrated circuit device, and particularly for types of integrated circuit devices, such as programmable devices, where a user may specify a clock rate.\nCertain types of integrated circuit devices allow users to specify various settings, such as clock rates. In particular, programmable devices, including, for example, programmable logic devices such as field-programmable gate arrays (FPGAs), may allow a user to specify a complete logic configuration, various portions of which may require different clock rates, none of which are known with any certainty at the time of device manufacture. Such devices have been manufactured with circuitry to allow various clock rates to be selected by the user, which may have resulted in overly complex clock networks, including many components that may never be used by a particular user.\nFor example, such devices may incorporate high-speed serial interfaces to accommodate high-speed (i.e., greater than 1 Gbps) serial I/O standards. Because there are multiple different standards, which may operate at multiple different rates, and because a user may elect to use more than one standard and/or rate, the ability to provide multiple reference clocks may be desirable. Heretofore, this has required the provision of multiple reference clock sources such as phase-locked loops (PLLs) or delay-locked loops (DLLs), with a clock network capable of routing a reference clock signal from any one of those sources to any one of a number of interface circuits."}
{"text": "Most specimens that are observed with a microscope have small variations in height across their surfaces. While these variations are frequently not visible to the human eye, they can cause images of a portion of a specimen captured by a microscope to be out of focus.\nThe range in which a microscope can create a usable focused image is known as the depth of field. The microscope must keep a portion of a specimen within its depth of field to generate useful images. However, when transitioning from observing a first portion of a specimen to observing a second portion of the specimen, the small variations in height of the specimen may cause the second portion to be outside the depth of field.\nDifferent sharpness measurements such as image contrast, resolution, entropy and/or spatial frequency content, among others, can be used to measure the quality of focus of images captured by a microscope. Generally, when a specimen is in focus, the captured image will exhibit the best sharpness quality (e.g., large contrast, a high range of intensity values and sharp edges). The different sharpness measurements that can be used to determine when a specimen is in focus usually require capturing a series of images and increasing or decreasing the distance between the microscope objective lens and the specimen until the image appears in focus. This increases the total microscopic scan time of each specimen, making methods using such measurement prohibitively slow for high throughput scanning applications.\nAccordingly, it is desirable to find a suitable in-focus plane of a specimen using a smaller number of images."}
{"text": "Anti-fuse is one of the One-Time Programmable (OTP) devices that can only be programmed once. Particularly, an anti-fuse has a high impedance state after fabrication and a low impedance state after being programmed. On the contrary, a fuse has a low impedance state after fabrication and a high impedance state after being programmed. The most commonly used anti-fuses are based on MOS gate oxide breakdown, metal-dielectric-metal breakdown, metal-dielectric-silicon breakdown, or silicon-dielectric-silicon breakdown, etc. Silicon dioxide (SiO2) is the most commonly used dielectric for breakdown in anti-fuses. However, Silicon-Oxide-Nitride (SON), Silicon Nitride (SiNx), Oxide-Nitride-Oxide (ONO), or other type of metal oxides, such as Aluminum Oxide (Al2O3), MgO, HfO2, or Cr2O3, can also be used.\nMOS gate oxide breakdown is based on applying a high voltage to break down the gate oxide to create a programmed state. However, there is a mechanism called soft-breakdown, other than the desirable hard-breakdown, which makes the dielectric film appear to be broken down, but the film may heal by itself after cycling or burn-in. The reliability may be a concern for practical applications.\nDielectric breakdown anti-fuses have been proven in manufacture. One of conventional dielectric breakdown anti-fuse is shown in FIGS. 1(a), 1(b), and 1(c). This anti-fuse is based on metal-dielectric-silicon with a diode constructed by P+ active region over N+ bar as program selector. FIG. 1(a) shows a portion of process steps by using a first Local Oxidation (LOCOS) to define an N+ bar area. FIG. 1(b) shows a second LOCOS step to further define active regions within each N+ bar in a perpendicular direction. The cell is patterned by two LOCOS steps so that the cell size is determined by the pitches of active regions in the X- and Y-directions. The cell size is generally referred to 4F2, where F stands for figure size. After the active region of the cells is determined, a P type dopant is implanted, a thin silicon dioxide is grown, and then a metal is built on top of each cell as shown in FIG. 1(c). The equivalent circuit of the anti-fuse cell is a capacitor in series with a diode at an X and Y cross-point as shown in FIG. 1(d). For additional information see, e.g., Noriaki, et. al, “A New Cell for High Capacity Mask ROM by the Double LOCOS Techniques,” International Electronics Device Meeting, December, 1983, pp. 581-584.\nThe anti-fuse cell in FIGS. 1(a), 1(b), and 1(c) is very complicated to fabricate, as it requires three more masks and two LOCOS steps over standard CMOS processes. Fabricating LOCOS requires a mask for field implant, nitride deposition, and a long thermal cycle to grow field oxide. Accordingly, there is a need for an anti-fuse cell that is more compatible with standard CMOS process to save costs."}
{"text": "Technical Field of the Invention\nThis invention relates generally to computing systems and more particularly to data storage solutions within such computing systems.\nDescription of Related Art\nComputers are known to communicate, process, and store data. Such computers range from wireless smart phones to data centers that support millions of web searches, stock trades, or on-line purchases every day. In general, a computing system generates data and/or manipulates data from one form into another. For instance, an image sensor of the computing system generates raw picture data and, using an image compression program (e.g., JPEG, MPEG, etc.), the computing system manipulates the raw picture data into a standardized compressed image.\nWith continued advances in processing speed and communication speed, computers are capable of processing real time multimedia data for applications ranging from simple voice communications to streaming high definition video. As such, general-purpose information appliances are replacing purpose-built communications devices (e.g., a telephone). For example, smart phones can support telephony communications but they are also capable of text messaging and accessing the internet to perform functions including email, web browsing, remote applications access, and media communications (e.g., telephony voice, image transfer, music files, video files, real time video streaming. etc.).\nEach type of computer is constructed and operates in accordance with one or more communication, processing, and storage standards. As a result of standardization and with advances in technology, more and more information content is being converted into digital formats. For example, more digital cameras are now being sold than film cameras, thus producing more digital pictures. As another example, web-based programming is becoming an alternative to over the air television broadcasts and/or cable broadcasts. As further examples, papers, books, video entertainment, home video, etc. are now being stored digitally, which increases the demand on the storage function of computers.\nA typical computer storage system includes one or more memory devices aligned with the needs of the various operational aspects of the computer's processing and communication functions. Generally, the immediacy of access dictates what type of memory device is used. For example, random access memory (RAM) memory can be accessed in any random order with a constant response time, thus it is typically used for cache memory and main memory. By contrast, memory device technologies that require physical movement such as magnetic disks, tapes, and optical discs, have a variable response time as the physical movement can take longer than the data transfer, thus they are typically used for secondary memory (e.g., hard drive, backup memory, etc.).\nA computer's storage system will be compliant with one or more computer storage standards that include, but are not limited to, network file system (NFS), flash file system (FFS), disk file system (DFS), small computer system interface (SCSI), internet small computer system interface (iSCSI), file transfer protocol (FTP), and web-based distributed authoring and versioning (WebDAV). These standards specify the data storage format (e.g., files, data objects, data blocks, directories, etc.) and interfacing between the computer's processing function and its storage system, which is a primary function of the computer's memory controller.\nDespite the standardization of the computer and its storage system, memory devices fail; especially commercial grade memory devices that utilize technologies incorporating physical movement (e.g., a disc drive). For example, it is fairly common for a disc drive to routinely suffer from bit level corruption and to completely fail after three years of use. One solution is to utilize a higher-grade disc drive, which adds significant cost to a computer.\nAnother solution is to utilize multiple levels of redundant disc drives to replicate the data into two or more copies. One such redundant drive approach is called redundant array of independent discs (RAID). In a RAID device, a RAID controller adds parity data to the original data before storing it across the array. The parity data is calculated from the original data such that the failure of a disc will not result in the loss of the original data. For example, RAID 5 uses three discs to protect data from the failure of a single disc. The parity data, and associated redundancy overhead data, reduces the storage capacity of three independent discs by one third (e.g., n−1=capacity). RAID 6 can recover from a loss of two discs and requires a minimum of four discs with a storage capacity of n−2.\nWhile RAID addresses the memory device failure issue, it is not without its own failure issues that affect its effectiveness, efficiency and security. For instance, as more discs are added to the array, the probability of a disc failure increases, which increases the demand for maintenance. For example, when a disc fails, it needs to be manually replaced before another disc fails and the data stored in the RAID device is lost. To reduce the risk of data loss, data on a RAID device is typically copied on to one or more other RAID devices. While this addresses the loss of data issue, it raises a security issue since multiple copies of data are available, which increases the chances of unauthorized access. Further, as the amount of data being stored grows, the overhead of RAID devices becomes a non-trivial efficiency issue."}
{"text": "1. Field of Invention\nThis invention relates to display units such as used by retail establishments for merchandising various wares. More particularly, this invention relates to a vertically extensible bar and a clutch mechanism for holding it in adjusted position.\n2. Description of the Prior Art\nA vertically extensible bar of a display rack may have an arm for supporting a series of hangers for clothes or other merchandise. Optionally, it may cooperate with a companion bar for supporting a shelf. The typical prior art structure for holding the bar in an extended position is a series of holes in the standard and a spring detent in the bar. There are several drawbacks to this arrangement. One disadvantage is that the series of holes in the standard are unsightly. Merchandisers appreciate more elegance in the display units for their merchandise.\nOther disadvantages include the inability to achieve infinite adjusted positions, the possibility of the coupling inadvertently slipping, the necessity of performing fabrication steps both on the standard and the bar."}
{"text": "1. Field of the Invention\nThe present disclosure generally relates to hinge assemblies and, particularly, to a hinge assembly used in a foldable electronic device having a top cover and a main body.\n2. Description of the Related Art\nFoldable electronic devices, such as notebook computers, are popular for their portability. In the foldable electronic device, a cover is rotatably connected to a main body via a typical hinge assembly.\nThe typical hinge assembly often includes a shaft, a friction member fixed on the shaft, and a rotary member rotatably sleeved on the shaft. The rotary member and the shaft are fixed to the cover and the main body. The rotary member is capable of being positioned in any position relative to the friction member and the shaft because of friction created between the friction member and the rotary member. Thus, the cover can be opened to any angle relative to the main body, and remain in any position.\nHowever, over time, the friction causes the engaging surfaces of the rotary member and the friction member to become abraded, resulting in little or no friction between the rotary member and the friction member. Thus, the cover would be incapable of remaining in any desired position. Therefore, the hinge assembly has a relatively short service life.\nTherefore, a new hinge assembly is desired to overcome the above-described shortcomings."}
{"text": "Golf is enjoyed by people of all ages throughout the world. It is played and enjoyed by both athletic types and non-athletic types of people. It is one of the few athletic activities that can be played and enjoyed by people with handicaps, such as a missing limb, a bad back, sight problems, and the like. The game is very challenging and requires mental concentration as well as athletic ability. Fortunately it does not require the athletic ability of a good or better athlete to play a decent round of golf. Through lessons and practice almost anyone can learn to play to a level that makes the game enjoyable and challenging. One of the greatest aspects of golf is that the golfer is playing against him or her self. A good golfer concentrates on his or her game, not the game or score of others he or she is playing with or against. This is a lesson that all great golf pros know.\nThere are typically four aspects to the game: the driving, the short game, getting out of trouble, and putting. Each aspect has its own challenges and requires practice to master, if ever mastered.\nPutting appears deceptively simple. It seems like anybody could grab a putter, i.e., the putting golf club, and putt or hit the golf ball to have the ball roll to the cup or hole and roll in it. Unfortunately putting is not simple and requires practice, more practice and even more practice. There is an old golf adage: golfers drive for show and putt for money. Great golf pros are great putters. They know the speed or condition of the putting green surface and their caddie, a professional just as much as the pro golfer, has read and recorded the green slope[s] from many different directions. This information is imparted to the pro by the caddie. It is very common to see the golf pro conferring with his or her caddie when the pro has to putt. The caddie frequently refers to a notebook that they prepare prior to the tournament regarding their readings of the putting green. These readings refer to the slope of the green surface or grass and the slope[s] of the green between the golf ball and the cup or hole.\nMost golfers cannot afford a full time caddie who reads the course for them before they begin a round of golf or who is familiar with the course from caddying on and playing on the course. At those golf courses and clubs that still make available experienced caddies, the golfer can learn about a putting green from the experienced caddie during a round of golf.\nThere are three conditions that must be taken into account when putting: (1) the rolling speed of the golf ball on the putting surface, (2) the slope[s] of the putting surface between the golf ball and the hole, and (3) the condition of the putting surface, such as, type and condition of the surface. For purposes of this patent, a course means a golf course or any practice facility, a ball means a golf ball, a putting surface means the putting green or practice surface, a cup or hole means the hole in the green which the golfer is attempting to get his or her ball into, a putter is the golf club that the golfer uses to hit a ball on the green in an attempt to roll the ball to and into the cup or hole, ‘to putt’ or ‘putting’ or ‘putt’ means to hit the ball on the green with the putter in an attempt to have the ball roll to and into the cup or hole, and ‘to sink a putt’ or ‘sinking a putt’ means to putt the ball and to successfully have the ball roll to and in the cup or hole. The direct putting path is the line of sight straight line between the hole and the ball. As most golfers know, a putt attempted along the direct putting path will in most instances, except for short putts, miss the hole for reasons set forth below. A short putt is normally a putt made three feet or less from the cup. The slope means a portion of the surface of a green that deviates from being parallel to the true level, true horizon or true level of the Earth (collectively “true level” herein). As used herein, putting green means putting surface. Most putting surfaces are closely cropped grass. However, synthetic putting surfaces, such as Astroturf® synthetic turf, packed oiled sand surfaces, dead grass surfaces, rug surfaces, exist.\nA spirit level is used to determine when an object, normally a surface of the object is parallel to true level or perpendicular to it. Some spirit levels are the arc of a circle and have a scale to give readings, in degrees, of the slope of the object to true level. To take, or to determine, or to measure a slope means to determine the angle of the surface of the green to the true horizon or true level of the Earth. The angle can be measure in degrees or it can be measured as an aiming point for the indicated putting path. Some slopes or breaks are perpendicular to the direct putting path; however most slopes or breaks intersect the direct putting path at angles other than 90 degrees. The actual putting path is the line, normally curved, between the ball and hole that will sink a putt when the ball is correctly hit or putted. The indicated putting path is the putting path indicated by the training device of the present invention. When the ball is properly putted, it will track the indicated putting path to sink a putt. The indicated putting path is determined by the measured slope and is the line of sight straight line between the ball and an aiming point on the training device.\nA golfer putts to the aiming point and because of the slope of the green and gravity the ball will normally roll along a curved path, the actual putting path. Because of the slope, the face of the putter must be square or perpendicular to the indicated putting path and gravity, the ball rolling along the indicated putting path will curve into the actual putting path which may have the ball curve away from the aiming point before the ball reaches the aiming point. No putt can be made if the golfer does not correctly hit the ball, i.e. impart the right amount of energy to the ball to roll to the cup, and address the putter face squarely to the aim point so the ball rolls, at least initially, towards the aiming point. The indicated putting path is normally a curved path between the ball and cup. To correctly putt, the face of the putter at impact must be square, i.e. perpendicular to the straight line intersecting the ball and the aim point\nIn most courses the hole is changed daily before the course is open to the public. In most courses the speed of each green, i.e. stimp rating, is taken daily and the stimp readings are available at the golf shop or pro shop. Where weather conditions can change dramatically during the course of the day, e.g. dewy or misty mornings and hot dry afternoons, the stimp readings may be taken two or three times a day. To successfully putt, a golfer must learn to adjust their putting or hitting of the golf ball to accommodate different green conditions which are indicative of different stimp readings. Every golfer has putted a ball to hard for the green conditions (a fast green-normally quiet dry) and has had the ball fly off the green and similarly every golfer has putted a ball too softly for the green conditions (a slow green-normally wet or moist grass, or dew) and had ball roll to only half as far as he she wanted or required to sink a putt.\nAlthough stimp readings are normally available for green conditions, readings on the lay of a green are not available. Greens are normally laid out and maintained with immaculate care. Some greens are as flat as a pancake and horizontal. If the ball is putted directly towards the hole on such greens, i.e. along the direct putting path, the ball will normally roll to and in the cup without deviation, assuming the ball was hit with the right amount of energy. However, most greens are sloped. Thus on a sloped green, even if the green is flat, gravity will effect the ball as it rolls placing a downhill force on the ball and the path of the ball will curve downhill. When putting on such greens, the golfer must compensate for the slope. Thus the putting path is a curved path on such greens. Unfortunately or fortunately, depending on the golfer's point of view, skill and experience, many greens have several slopes running in different directions. Thus as the hole is moved day to day and the ball most always lies in a new position on the green during each round, the golfer must learn to read the slope[s] of a green and adjust his or her putting to accommodate for the slope. As a rule, when putting at distances of less than twenty five feet, the golfer will only have a single slope to contend with. With longer putting distances, and even with shorter putting distances on some greens, the golfer may have contend with two or more slopes between the ball and the hold. These are very difficult putts and a lot of luck is required to make the putt.\nMost golfers after a little experience will attempt to read green slope[s]. Normally if the putt is a long one the golfer will walk the length of the distance between the ball or around the path, but observing golf etiquette-no walking on the path between the hole and another player's ball, and view the putting path both from the ball and the cup. The golfer will frequently squat down and attempt to read, or at least observe, the slope[s] of the green along the putting path. This can be difficult when the light is even, especially when the sun is overhead or the sky is clouded over, because shadowing is minimal. It is normally easier to read the slope[s] when the sun is low and the sun rays effect shadowing on the surface of green. As mentioned above, the professional caddie will have mapped out the slopes of the green before a tournament and an experience caddie at the courses that have them available, a rarity today, will have mentally mapped the slopes of the greens from experience on the course. However for the great majority of golfers in this world, the use of such caddie's is a rarity and the golfers must learn to read the slopes on their own. It even takes golfers who play the same course week end and week out at least six months to memorize the breaks on the course greens.\nThe object of the present invention is to provide a golfer with a device that can measure the slope of a particular area of a green. The training device can be used to confirm a golfer's slope reading of a particular area of the green.\nA further object of the present invention is to provide a means for a golfer to read an indicated slope of a particular area of the green and then confirm the reading by measuring the slope with the device of the present invention.\nIt is a further object of the present invention to provide a device that can give the golfer an aim point for putting to give the indicated putting path so the golfer learns to determine the aim point upon determining the slope. Most golfers will not determine slope in degrees. Most golfers will develop an intuition about the slope and correlate the aim point to the slope. With experience, the golfer will correlate the speed or slowness of the surface and the conditions of the putting surface into his or her determination of the aim point.\nA still further object of the present invention is provide a means for a golfer to read the slope of a particular area of the green, to determine the indicated putting path based on the slope reading, and then confirm the slope reading by measuring the slope and confirm the indicated putting path by attempting the putt along the indicated putting path.\nA golfer can meet the present objectives by employing the device on a practice putting green to learn how to read green slopes, and/or to determine the indicated putting path, and then confirming slope readings with the device and/or confirming the indicated putting path by putting. A golfer can also use the device during a game assuming his or her partners don't object."}
{"text": "1. Technical Field\nThe disclosure relates in general to a memory device and a manufacturing method thereof, and particularly to a memory device having a reduced size as well as an excellent operating performance and a manufacturing method thereof.\n2. Description of the Related Art\nConventionally, in a manufacturing process for forming a memory device, a whole polysilicon film is deposited and then etched to form word lines. Next, dielectric materials are filled into the spaces between the word lines. However, as the reduction of the sizes of memory devices, the widths of word lines and between which the gaps are reduced as well. As such, the word lines may be short-circuited due to the residual polysilicon between the word lines manufactured by etching processes, caused by an incomplete etching between the word lines, or the widths of word lines are not uniform, resulting in lower reliability of the memory devices. In addition, the reduction of widths of word lines results in poor performances of memory devices.\nAccordingly, it is desirable to develop memory devices with improved reliability and operating performance."}
{"text": "Many articles of furniture are costly to ship because they are by nature bulky and prone to damage during transport. Therefore, it has been common to make knock down type mass-market furniture. Knock down furniture is fabricated as components, or sub-assemblies, which can be compactly packaged and economically shipped. The furniture is subsequently assembled by a retailer or a consumer using simple tools, such as common wrenches, screwdrivers, hexagonal wrenches, and the like. Most often such furniture can be subsequently disassembled, if desired. However, the advantages of knock down design will not be realized if such a design compromises the article's appearance or function, or if the article is too hard to assemble.\nWhat constitutes a compromise in appearance for a knock down article depends on an esthetic judgment, and that may vary with the individual. Nonetheless, there are some general principles which may be stated. For example, most people would conclude it is esthetically undesirable to have exposed industrial-type metal fasteners on a wooden chair. Similarly, if the knock down design involved significant changes in the proportions or shapes of the parts of a chair, compared to a traditional chair design which was obviously being emulated, then there would be a high risk that consumers would think the chair looked strange, and they would not purchase it.\nA knock down design which compromises function becomes evident when the piece of furniture is put into use. A chair may be subjected to very high loads. For instance, the chair may set on an uneven surface, a user may tilt the chair backward on the rear legs, or the chair may fall over onto a hard floor. Consequently, a knock down chair must not only have strength and rigidity when first assembled, but it must maintain such during its lifetime.\nIn furniture which is factory-assembled, it is possible to use heavy machinery and special processes. It is possible to use tight fits, diverse fasteners, and special adhesives; all to obtain the strength and durability the product demands. In contrast, by the nature of knock down furniture, there will be joints which must be made by unskilled consumers using simple hand tools. Thus, in some poorly designed knock down furniture the joints will be weak and furniture will be flimsy when initially assembled. In other such furniture, joints will loosen with time or even fail during use. In still other furniture, the knock down design may provide good strength, but be too complex for unskilled consumers to assemble correctly. And of course, a piece of knock down furniture has to be economic to manufacture, otherwise the advantage produced by lower packaging and transport costs, compared to a one-piece factory assembled chair, will be offset.\nSo, it is not easy to make a piece of knock down furniture which satisfactorily meets all the requirements. Of course, there have been many successful designs of knock down furniture. Specialized fasteners have been developed. However, certain designs of furniture by their nature still present problems which are more difficult to overcome than others. For example, joints which are made at obvious locations can be subject to inherently high stresses, as is the case when a cantilevered back rest of a chair is joined to the chair seat. Therefore, there is a continuing search for new knock down concepts and joint designs."}
{"text": "The present invention is directed to a three dimensional (3-D) graphics application programming interface (API) that provides new and improved methods and techniques for communications between application developers and procedural shaders, such as vertex and pixel shaders.\nComputer systems are commonly used for displaying graphical objects on a display screen. The purpose of three dimensional (3-D) computer graphics is to create a two-dimensional (2-D) image on a computer screen that realistically represents an object or objects in three dimensions. In the real world, objects occupy three dimensions. They have a real height, a real width and a real depth. A photograph is an example of a 2-D representation of a 3-D space. 3-D computer graphics are like a photograph in that they represent a 3-D world on the 2-D space of a computer screen.\nImages created with 3-D computer graphics are used in a wide range of applications from video entertainment games to aircraft flight simulators, to portray in a realistic manner an individual\"\"s view of a scene at a given point in time. Well-known examples of 3-D computer graphics include special effects in Hollywood films such as Terminator II, Jurassic Park, Toy Story and the like.\nOne industry that has seen a particularly tremendous amount of growth in the last few years is the computer game industry. The current generation of computer games is moving to 3-D graphics in an ever increasing fashion. At the same time, the speed of play is driven faster and faster. This combination has fueled a genuine need for the rapid rendering of 3-D graphics in relatively inexpensive systems.\nRendering and displaying 3-D graphics typically involves many calculations and computations. For example, to render a 3-D object, a set of coordinate points or vertices that define the object to be rendered must be formed. Vertices can be joined to form polygons that define the surface of the object to be rendered and displayed. Once the vertices that define an object are formed, the vertices must be transformed from an object or model frame of reference to a world frame of reference and finally to 2-D coordinates that can be displayed on a flat display device, such as a monitor. Along the way, vertices may be rotated, scaled, eliminated or clipped because they fall outside of a viewable area, lit by various lighting schemes and sources, colorized, and so forth. The processes involved in rendering and displaying a 3-D object can be computationally intensive and may involve a large number of vertices.\nTo create a 3-D computer graphical representation, the first step is to represent the objects to be depicted as mathematical models within the computer. 3-D models are made up of geometric points within a coordinate system consisting of an x, y and z axis; these axes correspond to width, height, and depth respectively. Objects are defined by a series of points, called vertices. The location of a point, or vertex, is defined by its x, y and z coordinates. When three or more of these points are connected, a polygon is formed. The simplest polygon is a triangle.\n3-D shapes are created by connecting a number of 2-D polygons. Curved surfaces are represented by connecting many small polygons. The view of a 3-D shape composed of polygon outlines is called a wire frame view. In sum, the computer creates 3-D objects by connecting a number of 2-D polygons. Before the 3-D object is ultimately rendered on a 2-D display screen, however, the data of sophisticated graphics objects undergoes many different mathematical transformations that implicate considerably specialized equations and processing unique to 3-D representation.\nAs early as the 1970s, 3-D rendering systems were able to describe the xe2x80x9cappearancexe2x80x9d of objects according to parameters. These and later methods provide for the parameterization of the perceived color of an object based on the position and orientation of its surface and the light sources illuminating it. In so doing, the appearance of the object is calculated therefrom. Parameters further include values such as diffuse color, the specular reflection coefficient, the specular color, the reflectivity, and the transparency of the material of the object. Such parameters are globally referred to as the shading parameters of the object.\nEarly systems could only ascribe a single value to shading parameters and hence they remained constant and uniform across the entire surface of the object. Later systems allowed for the use of non-uniform parameters (transparency for instance) which might have different values over different parts of the object. Two prominent and distinct techniques have been used to describe the values taken by these non-uniform parameters on the various parts of the object\"\"s surface: procedural shading and texture mapping. Texture mapping is pixel based and resolution dependent.\nProcedural shading describes the appearance of a material at any point of a 1-D, 2-D or 3-D space by defining a function (often called the procedural shader) in this space into shading parameter space. The object is xe2x80x9cimmersedxe2x80x9d in the original 1-D, 2-D or 3-D space and the values of the shading parameters at a given point of the surface of the object are defined as a result of the procedural shading function at this point. For instance, procedural shaders that approximate appearance of wood, marble or other natural materials have been developed and can be found in the literature.\nThe rendering of graphics data in a computer system is a collection of resource intensive processes. The process of shading i.e., the process of performing complex techniques upon set(s) of specialized graphics data structures, used to determine values for certain primitives, such as color, etc. associated with the graphics data structures, exemplifies such a computation intensive and complex process. For each application developer to design these shading techniques for each program developed and/or to design each program for potentially varying third party graphics hardware would be a Herculean task, and would produce much inconsistency.\nConsequently, generally the process of shading has been normalized to some degree. By passing source code designed to work with a shader into an application, a shader becomes an object that the application may create/utilize in order to facilitate the efficient drawing of complex video graphics. Vertex shaders and pixel shaders are examples of such shaders.\nPrior to their current implementation in specialized hardware chips, vertex and pixel shaders were sometimes implemented wholly or mostly as software code, and sometimes implemented as a combination of more rigid pieces of hardware with software for controlling the hardware. These implementations frequently contained a CPU or emulated the existence of one using the system\"\"s CPU. For example, the hardware implementations directly integrated a CPU chip into their design to perform the processing functionality required of shading tasks. While a CPU adds a lot of flexibility to the shading process because of the range of functionality that a standard processing chip offers, the incorporation of a CPU adds overhead to the specialized shading process. Without today\"\"s hardware state of the art, however, there was little choice.\nToday, though, existing advances in hardware technology have facilitated the ability to move functionality previously implemented in software into specialized hardware. As a result, today\"\"s pixel and vertex shaders are implemented as specialized and programmable hardware chips. Exemplary hardware designs of vertex and pixel shader chips are shown in FIGS. 1A and 1B, and are described later in more detail. These vertex and pixel shader chips are highly specialized and thus do not behave as CPU hardware implementations of the past did.\nThus, a need has arisen for a 3-D graphics API that exposes the specialized functionality of today\"\"s vertex and pixel shaders. In particular, since present vertex shaders are being implemented with a previously unheard of one hundred registers, it would be advantageous to have a register index for indexing the registers of the vertex shader. Also, since realistic simulations require the precision of floating point numbers, it would be advantageous to provide specialized vertex shading functionality with respect to the floating point numbers at a register level. For example, it would be desirable to implement an instruction set that causes the extremely fast vertex shader to return only the fractional portion of floating point numbers. Similarly, with respect to pixel shaders, it would be desirable to provide specialized pixel shading functionality as well. More particularly, it would be desirable to provide a function that performs a linear interpolation mechanism. Furthermore, it would be desirable to use operation modifiers in connection with an instruction set tailored to pixel shaders. For example, negating, remapping, biasing, and other functionality would be extremely useful for many graphics applications for which efficient pixel shading is desirable, yet as they are executed as part of a single instruction they are best expressed as modifiers to that instruction. In short, the above functionality would be advantageous for a lot of graphics operations, and their functional incorporation into already specialized pixel and vertex shader sets of instructions would add tremendous value from the perspective of ease of development and improved performance.\nIn view of the foregoing, the present invention provides a three-dimensional (3-D) API for communicating with hardware implementations of vertex shaders and pixel shaders having local registers. With respect to vertex shaders, API communications are provided that may make use of an on-chip register index and API communications are also provided for a specialized function, implemented on-chip at a register level, which outputs the fractional portion(s) of input(s). With respect to pixel shaders, API communications are provided for a specialized function, implemented on-chip at a register level, that performs a linear interpolation function and API communications are provided for specialized modifiers, also implemented on-chip at a register level, that perform modification functions including negating, complementing, remapping, biasing, scaling and saturating. Advantageously, the API communications of the present invention expose very useful on-chip graphical algorithmic elements to a developer while hiding the details of the operation of the vertex shader and pixel shader chips from the developer.\nOther features of the present invention are described below."}
{"text": "The present invention relates to a system of locating and identifying underground pipes, such as those which carry gas, water and waste to/from homes and building so that, among other things, these pipes can be avoided by excavating equipment or the like. More particularly, the system of the present invention enables the detection of such objects utilyzing, in combination with other novel elements of this invention, ground-probing RADAR.\nThis invention further relates to the RADAR detection of underground pipes, for example, and more particularly relates to a system which focuses a characteristic hyperbolic RADAR response received from an underground object, such as a pipe or other object, utilizing synthetic aperture type technologies, and further processes same to accurately determine said object's underground position.\nAccurate RADAR-based underground object detection has always been an elusive goal because of the variability of the ground as a conducting medium in three dimensions, i.e., inherent variations in ground layers, density, obstructions, dielectric constant, etc. Water content, in particular, acutely varies the ground's dielectric constant which correspondingly attenuates RADAR signals making consistent detection of targets underground difficult at best. Electromagnetic signals transmitted into the ground and reflected from an object buried therein tend to suffer high signal attenuation resulting in low signal-to-clutter and signal-to-noise ratio. Efforts to improve detection ability have found that while a single frequency of operation may be desirable in a particular type soil, the same frequency may be undesirable in another, frequently misinterpreting said objects as ground clutter by conventional underground radar systems.\nIn an effort to overcome inadequacies of conventional underground RADAR detection, U.S. Pat. No.3,831,173 discloses a ground radar system which utilizes a transient signal comprising a wide variety of radiated frequencies, Due to the use of the transient signal, effective reflections are received from a wide variety of underground objects such as pipes, utility lines, culverts, ledges, etc., to depths around 10 feet. The '173 system, however, while appropriate for detecting small conducting objects, is basically unable to accurately detect non-conducting objects with cross-sections of less than one or two feet.\nU.S. Pat. No. 3,967,282 discloses a detector for detecting both metallic and non-metallic objects, based on differences in the dielectric constants of the object and its surrounding medium, in order to give a location of the object. The '282 invention, however, is burdened with difficulty in processing the received data such that accurate object detection and positioning is not achieved.\nU.S. Pat. No. 4,706,031 discloses a method and apparatus for identifying a target object located in the ground, in the air or under water. The basis within the disclosure for detection and target identification resides in the apparatus use of the phase deviation between the transmitted and received (echo) radio-wave signals. A signal containing a mixture of various frequencies is transmitted and the return signal or signals are analyzed. A detected difference in phase deviation between the particular frequencies received is used to identify the material properties of the object from which the energy is reflected. The '031 apparatus, however, is still plagued with problems when it comes to detecting small non-conducting objects.\nU.S. Pat. No. 4,951,055 discloses a ground probing RADAR which includes means for displaying echo images of a buried material. The displayed images are capable of providing a depth direction of the buried material and a movement direction of a moving vehicle carrying the RADAR. The RADAR includes first means for forming a hyperbolic echo image of the material, and causing a hyperbolic echo image to be displayed on the display means, second means for forming a false echo image and causing the false echo image to be displayed on the display means, third means for inputting data to the second means to cause a displayed position of the false echo image to be shifted so that a vertex position of the false echo image and expansion opening thereof coincide with those of the echo image of the buried material, and fourth means for calculating a propagation velocity of the electromagnetic waves in the ground on the data indicative of the vertex position and opening expansion of the false echo image when the two displayed echo images coincide with each other. A position of the buried material is detected on the basis of the propagation velocity value calculated by the fourth means.\nThat is, when electromagnetic waves are emitted from a plurality of points on the ground surface above a buried material, an echo image formed on the basis of data of propagation times of reflected waves at their respective points describes a hyperbola as a result of expansion of the transmitted electromagnetic waves. An operation is carried out to overlap, on the echo image, a false echo image lying in the same coordinate system and consisting of a similar hyperbolic image. If the two echo images are overlapped, a vertex position and an expansion of the opening of the echo image can be determined from the data of the false echo image. Thus, the propagation velocity of electromagnetic waves are calculated from the data that represents the vertex position and the expansion of the opening. The position of the material under the ground is then calculated in relation to the data of propagation time in any position."}
{"text": "The present disclosure relates to an image processing apparatus and a log management method for storing logs of jobs into a storage portion.\nIn image processing apparatuses such as printers, scanners, or the like, logs of various types of jobs, such as a print job, a scan job, or the like, are stored and accumulated in a storage portion. The logs may include image logs generated based on image data which is the processing target of the job. In this type of image processing apparatus, it is known that, when the storage capacity of the storage portion in which the logs are stored has become equal to or lower than a predetermined amount, logs are deleted in order from the oldest log."}
{"text": "The majority of the papermaking pulp produced in the world today is produced by the so-called kraft method. Kraft pulping produces strong fibers, a fact that has given the method its name. This method, however, has the drawback of being very capital intensive. This is due to the need for a very complex system for chemicals recovery and very large unit sizes in the reactors. The reactors have in fact become so big that controlling the actual reactions and liquor circulations has become extremely difficult. The huge unit sizes in all parts of the process also leads to very large in-process inventory and a process that reacts very slowly to e.g. grade changes, etc. Any improvement that would lead to a faster process with shorter in-process delays would therefore have to be seen as a big step forward.\nAnother problem regarding the kraft method is the use of sulfur, which leads to larger amounts of chemicals being in circulation, odor problems, as well as making the recovery of spent chemicals extra complicated. A process without sulfur would make it possible to have much more efficient burning processes for the dissolved organic material in the process.\nIn order to address the problems of slow and cumbersome processes and to get rid of the sulfur, and often all inorganic chemicals in the process, several researchers have proposed the use of organic solvents to act as a cooking chemical and dissolve the lignin that holds the cellulose fibers together in wood.\nAccording to J. Gullichsen, C-J Fogelholm, Book 6A, Papermaking Science and Technology, Fapet, 1999, Helsinki, Finland, p. B411, the pulping methods using organic solvents can be classified as follows:                Autohydrolysis methods, in which organic acids released from the wood by thermal treatment act as delignification agents        Acid catalyzed methods, in which acid agents are added to the material        Methods using phenols        Alkaline organosolv methods        Sulfite and sulfide cooking in organic solvents        Cooking using oxidation of lignin in organic solvent        \nThe basic idea in autohydrolysis, as explained for instance in U.S. Pat. No. 3,585,104 (Kleinert), is to cook the wood in a solvent at high temperature. The high temperature leads to hydrolysis of sugars present in the wood, thus releasing acids. These acids are then supposed to break down and dissolve lignin together with the solvent. The drawback of this process is that very harsh conditions are needed in order to properly delignify the wood. This leads to yield losses and low pulp quality. Others have attempted to improve on the basic idea in order to improve the pulp quality. One such attempt is the so-called IDE process described in EP 0 635 080. The idea is to limit the drop in pH in order to salvage pulp quality. The process is proposed to achieve this by cooking using solvent in a countercurrent manner, thus removing the acids as they are formed early in the cook, and by adding alkali to maintain the pH as desired. The method has never been possible to implement on a commercial scale, possibly due to the large amount of solvent needed to maintain the proposed countercurrent flow. Further, even in the laboratory it is not well suited for all wood species.\nIf pulp quality is not seen as a major criteria (emphasis on by-product value), acid can be added to the system to increase the speed of the pulping process. Processes have for instance been developed that use acetic and formic acid as delignification agents. The drawback for these processes is that there is no market for the inferior quality pulp, and that severe corrosion problems arise in the equipment.\nThe so-called Organocell process has been closest to large-scale commercialization of the solvent-using pulping methods. This process is a variant of alkaline organosolv pulping, using simultaneous action of soda-anthraquinone and organic solvent on the lignin. The process seemed to give acceptable pulp quality in the laboratory, but when tried on mill scale the results were not satisfactory.\nAll prior pulping methods employing organic solvents have been attempts to develop substitutes for the presently dominating kraft pulping method. However, kraft pulping has been constantly improved upon for the last 100 years and is today quite efficient and thus hard to compete with. This can be seen from the fact that no solvent pulping method has proven to be commercially viable. There is, however, still room for improvement in the kraft process itself. For example, the odors of the process are seen as a problem, as is the fact that the reactors are becoming increasingly large and hard to control. Steps have been taken to improve alkaline kraft pulping. One such method is rapid steam phase pulping. The idea is to impregnate the wood with all the alkaline chemicals needed for the reactions in an impregnation stage, followed by heating in a water steam phase. This would make the reactors smaller and partly remedy the problems with odor as described in Canadian Patent No. 725,072. However, this method has not demonstrated enough improvement over the kraft process in liquid phase—yield increase has been very small and reactors still very big, leading to too high chip columns in vapor phase, in turn leading to compaction and collapsing of the digester content, thus plugging flows and destroying pulp quality.\nIn light of the current research it is clear that the previous research has failed largely because the true role of the organic solvent was not identified. In the current research it has been clearly seen that organic solvents do not participate in the reactions themselves as a solvent of lignin or active chemical, but in fact only have the impact of providing such a reaction environment as to boost the efficiency of other delignifying chemicals."}
{"text": "1. Field of the Invention\nThe present invention relates to an anti-reflection film which is arranged for the purpose that external light is prevented from the reflection on the surface of a window, a display, and the like. Particularly, the present invention relates to an anti-reflection film which is arranged on the surface of a display such as a liquid crystal display (LCD), a CRT display, an organic electroluminescence display (ELD), a plasma display (PDP), a surface-conduction electron-emitter display (SED), and a field-emission display (FED).\nParticularly, the present invention relates to an anti-reflection film which is arranged on the surface of a liquid crystal display (LCD). Further, the present invention relates to an anti-reflection film which is arranged on the surface of a transmission type liquid crystal display (LCD).\n2. Description of the Related Art\nIn general, displays are used in the environment into which external light and the like enter regardless of whether displays are used in indoor or outdoor. The incident light such as the external light causes regular reflection on the display surface and the like so that the reflected image is mixed with the displayed image and the quality of display screen is reduced. Therefore, it is essential to provide a display surface and the like with an anti-reflection function, and further, improvements of the anti-reflection function as well as a complex of functions other than the anti-reflection function are being demanded.\nIn general, an anti-reflection function is obtained as a result of the formation of an anti-reflection layer with a multilayer structure which repeatedly has high refractive index layers and low refractive index layers that are made of a transparent material such as metal oxide on a transparent substrate. The anti-reflection layer composed of the multilayer structure can be formed by a dry film-forming method such as chemical vapor deposition (CVD) and physical vapor deposition (PVD).\nIn the case where an anti-reflection layer is formed by a dry film-forming method, while there is an advantage of finely controlling the thickness of a low refractive index layer and a high refractive index layer, there is a problem of low productivity since the film is formed in a vacuum, which is thus unsuitable for mass production. On the other hand, as a method of forming an anti-reflection layer, the production of anti-reflection film by a wet film-forming method with the use of a coating liquid in which a large area, continuous production, and cost reduction are possible, has been attracting attention.\nIn addition, in an anti-reflection film in which such anti-reflection layer is arranged on a transparent substrate, the surface is relatively flexible, therefore, in order to give hardness to the surface, a technique in which a hard coat layer that is obtained by curing of an acrylic-based material is arranged and an anti-reflection layer is formed on the hard coat layer is generally used. This hard coat layer is provided with a high level of surface hardness, luster, transparency, and excoriation resistance by the acrylic-based material.\nIn the case where an anti-reflection layer is formed by a wet film-forming method, the anti-reflection layer is produced with the application of at least a low refractive index layer on a hard coat layer that is obtained by curing of the ionizing radiation curable materials, and the wet film-forming method has a merit of inexpensive production in comparison with a dry film-forming method, and thus, anti-reflection layers produced by such a wet film-forming method are widely distributed in the market.\n<Patent document 1>: JP-A-2005-202389.\n<Patent document 2>: JP-A-2005-199707.\n<Patent document 3>: JP-A-H11-92750.\n<Patent document 4>: JP-A-2007-121993.\n<Patent document 5>: JP-A-2005-144849.\n<Patent document 6>: JP-A-2006-159415.\n<Patent document 7>: JP-A-2007-332181.\nWith the arrangement of an anti-reflection film onto a display surface, by the anti-reflection function of the anti-reflection film, the reflection of external light can be suppressed to improve the contrast of the display in a bright place. Further, the transmittance can be improved at the same time, therefore, an image can be displayed brighter than usual. In addition, the anti-reflection film is also expected to have an energy-saving effect that can suppress the power consumption of the backlight, and the like.\nAs for an anti-reflection film, an anti-reflection film with low production costs is demanded. Further, an anti-reflection film having excellent anti-reflection performance or excellent optical properties free from interference irregularity is demanded. In addition, an anti-reflection film with high excoriation resistance is demanded since the anti-reflection film is arranged on a display surface. In addition, an anti-reflection film having antistatic function for the prevention of dust adhesion is demanded. In the present invention, the problem to be solved is to provide an anti-reflection film with low production costs, having excellent optical properties and excellent excoriation resistance and antistatic function."}
{"text": "Axle driving units incorporating hydraulic stepless transmissions have been used to drive the axles of self-propelled vehicles for many years. Generally such units include a hydraulic pump driven by an input shaft and a hydraulic motor having an output shaft drivingly connected through a differential to a pair of oppositely disposed axles. An example of such a unit is disclosed in U.S. Pat. No. 4,914,907. However, certain self-propelled vehicles perform tasks which require tight turning capabilities and conventional hydraulic transmissions which drive a pair of axles through a differential gear assembly are not particularly suited for such purposes. Instead, vehicles have been provided with axles which are independently driven by separate axle drive units such that turns are accomplished by rotating drive wheels on opposite sides of the vehicle at different speeds and/or in different directions. Further, certain such axle driving units for independently driving single axle have incorporated hydraulic transmissions. However, such axle driving units have required housings which are of substantial height and substantial width in order to accommodate the hydraulic pump and motor and the other necessary components. Accordingly, vehicles have required large body frames in order to accommodate two such axle driving units in a side-by-side disposition, thus ruling out use of the units on many small vehicles. Further, even where a large body frame is provided, the center of gravity of the vehicle tends to be higher than is desirable for good roadability due to the height of the axle driving units and the need to dispose the prime mover of the vehicle in an elevated position to efficiently drive the units. For example, in U.S. Pat. No. 5,127,215 a dual hydrostatic drive walk-behind mower is disclosed, but it can be readily seen that the axle driving units of this mower require substantial vertical and lateral space such that a large body frame is required. It will also be noted that due to the height of the transmission housings, the engine must be disposed in an elevated position which results in the vehicle having an undesirably high center of gravity. Moreover, multiple driving belts are required to drive the input shafts of the axle driving units. (See also, U.S. Pat. Nos. 4,809,796 and 5,078,222). In U.S. Pat. No. 4,819,508, a transmission system for working vehicles is disclosed which partially solves the problem of an undesirable center of gravity by reorienting the engine such that the crank shaft is horizontally disposed. However, the axle driving mechanism still occupies substantial vertical space on the body frame, making the center of gravity undesirably high. Further, reorientation of the engine complicates the drive belt systems for driving both the axle driving units and the mower blades.\nTherefore, it is an object of the present invention to provide an axle driving apparatus for independently driving axles on opposite sides of a vehicle.\nIt is another object of the present invention to provide an axle driving apparatus which includes side-by-side axle drive units incorporating hydraulic transmissions which require limited vertical or lateral space such that the axle driving apparatus can be used by small self-propelled vehicles, and such that vehicles utilizing such axle driving apparatus define low centers of gravity for improved roadability.\nYet another object of the present invention is to provide an axle driving apparatus having input shafts and a drive belt system which facilitates drivingly connecting the apparatus to the prime mover of the vehicle.\nStill another object of the present invention is to provide an axle driving apparatus which is inexpensive to manufacture and maintain."}
{"text": "Voltage sampling is traditionally used for analog-to-digital (A/D) conversion. In a voltage sampler, a sampling switch is placed between a signal source and a capacitor. Between two sampling moments, the capacitor voltage tracks the signal voltage accurately. At the sampling moment, the switch is turned off to hold the capacitor voltage. The two processes become increasingly difficult when the signal frequency increases. For a given accuracy, thermal noise and switching noise set a minimum allowable capacitance while the tracking speed set a maximum allowable capacitance or switch resistance. It becomes impossible when the maximum is smaller than the minimum. Moreover, the clock jitter and finite turning-off speed (nonzero sampling aperture) make the sampling timing inaccurate. In fact, the bandwidth of a voltage sampling circuit must be much larger than the signal bandwidth. This makes direct sampling of high frequency radio signal extremely difficult. Sub-sampling can reduce the sampling rate but not the bandwidth of the sampling circuit and not the demands on small clock jitter and small sampling aperture."}
{"text": "Flame safeguard systems that utilize a flame sensor and amplifiers for control of valve means in a burner system have been utilized for many years. Typically these systems use discrete component electronic systems in their amplifiers, and the amplifiers in turn ultimately control a relay. The relay in turn is used to switch power for an electromagnetically operated fuel valve. As the electronic and electromechanical types of flame safeguard systems evolved, the reliability and safety of the systems has been of prime concern. As a result of this concern, systems and equipment have been developed which are very reliable and allow the flame safeguard system to accurately and reliably control the operation of the main fuel valve to a burner in response to the presence or absence of flame at the burner.\nIn recent years microcomputer based systems have evolved. These systems utilize very small and complex integrated circuits. The microcomputers, while having many capabilities, have a frailty in that they are subject to many more types of failures than discrete component electronic circuits. The utilization of a microcomputer in a flame safeguard control system requires a high degree of care, and the use of special safety systems. In a prior art type of microcomputer based flame safeguard system as disclosed in U.S. Pat. No. 4,298,334, assigned to the assignee of the present invention, the microcomputer and flame amplifier both controlled power to the main fuel valve relay. This type of redundant circuitry is expensive, and may be improved upon by the present invention."}
{"text": "The invention relates to a high-pressure fan which comprises an electric motor provided with a shaft that goes through it; blade wheels which are mainly made of a carbon fibre-based composite material, arranged on both sides of the electric motor and mounted directly on the shaft of the electric motor; fan housings surrounding the blade wheels; and an intermediate channel which connects a pressure opening of one fan housing to a suction opening of the other fan housing, the intermediate channel being integrated into the base structure of the high-pressure fan and located substantially below the electric motor and the blade wheels mounted to it.\nThis type of compact high-pressure fan except for the intermediate channel structure integrated to the base structure is known from FI patent 101564. When the blade wheels are made of a mainly carbon fibre-based composite material, the weight of the blade wheel is considerably reduced in comparison with the conventional steel and aluminium blade wheels. In spite of this, it has been possible to design the blade wheels at least as strong as when using the conventional materials. Due to the light weight of the blade wheels, they also do not require separate bearings and the bearings of the electric motor suffice. Due to the light weight of the rotating parts and the possibility to use smaller slide bearings instead of roller bearings, it has also been possible to raise considerably the rotation rate of the fan. The first critical rotation rate of the high-pressure fan is at the same time also at a higher level than that of the conventional solutions, which facilitates the use of rotating control in the entire area of operation, the control being easy to implement by means of a frequency converter.\nThe two series-connected steps of the high-pressure fan, i.e. the entities formed by each blade wheel and fan housing, are in this solution according to FI patent 101564 connected to each other by means of an intermediate channel running over the entire fan structure.\nEven though this known high-pressure fan structure provides the significant advantages listed above, problems are still caused by said intermediate channel that requires a lot of space and causes vibration especially at high rotation rates.\nAn intermediate channel structure integrated to the base structure is in turn known from SU invention report 1710849."}
{"text": "The invention relates generally to dielectric materials, and, in particular, to rare-earth-element-based titanate systems.\nA dielectric material is an insulating material that does not conduct electrons easily and thus has the ability to store electrical energy when a potential difference exists across it. Common dielectric materials include glass, mica, mineral oil, paper, paraffin, polystyrene, plastics, phenolics, epoxies, aramids, and porcelain. In electronic circuits, dielectric materials may be employed as capacitors. High dielectric constant materials may be used in radar or microwave applications and for circuit miniaturization as the speed of propagation of signal is related to the dielectric constant of the medium through which it passes. If the loss tangent for a material of a given frequency signal is very low, the electrical loss related to the hysteresis decreases resulting in an efficient signal transmission.\nThere is a need for a dielectric material that has one or more desirable characteristics, such as, a high dielectric constant, a low loss tangent, the ability to withstand a wide range of temperatures, the ability to operate in wide range of frequencies, voltages, atmospheric conditions, and pressures, and the capability for use in the manufacture of composite structures that can be used alone or in combination with other materials."}
{"text": "The present invention relates to a green concrete saw, or a saw for use in cutting grooves into wet concrete. The saw includes a base onto which a circular saw blade and engine are mounted, with a control handle that allows an operator control over the moving direction of the saw. The circular saw blade is rotated by the engine so as to cut grooves along a path of movement of the saw in the green or wet concrete. Accordingly, the base generally includes wheels that allow for easier movement of the saw.\nHowever, conventional green concrete saws are generally limited in their ability to adjust the height of circular saw blade relative to the ground or the wet concrete. As such, it is difficult to control the quality of the grooves cut into the wet concrete. Thus, there is a need for a green concrete saw with improved precision in setting the elevation of the circular saw blade relative to the wet concrete."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a method of fabricating a Metal Oxide Semiconductor Field Effect Transistor (MOSFET), and more particularly to a method of forming a field oxide film which provides hyperfine device isolation on a Silicon-on-Insulator (SOI) substrate by means of Local Oxidation of Silicon (LOCOS).\n2. Description of the Related Art\nWith the recent remarkable progress in semiconductor devices, demand is increasing for an LSI on which both digital and analog circuits are mounted, and which performs at high speed and with reduced power consumption. To meet this demand, semiconductor devices are required to be integrated more densely. As the devices to be mounted increase in number, isolation regions must be narrower and smaller.\nA conventional method of fabricating a MOSFET in an SOI substrate by means of LOCOS is illustrated in FIGS. 2A-2F, each of which schematically shows a cross-section of the MOSFET at a fabrication step. Descriptions of the steps are as follows:\na) A pad oxide film 52 of about 5-10 nm is deposited on an SOI substrate 51. Then an active nitride film 53 of about 50-150 nm is deposited on the pad oxide film 52 as an oxidation-resistant mask (see FIG. 2A).\nb) Openings are formed in the laminated layers of the pad oxide film 52 and the active nitride film 53 at positions where field oxide films 54 are to be provided, by a conventional lithography technique (see FIG. 2B).\nc) The field oxide films 54 are formed on the SOI substrate 51 by dry oxidation (a heat treatment conducted in a dry oxygen atmosphere) (see FIG. 2C).\nd) The remaining portions of the active nitride film 53 and the pad oxide film 52 are removed (see FIG. 2D).\ne) Gate electrodes 55 are provided by a conventional process for fabricating MOSFETs (see FIG. 2E).\nf) SiO2 side walls 57 are formed by first providing an SiO2 film on the substrate and then etching back. Impurities are then introduced into the substrate by means of ion implantation to form source/drain regions 58. Finally, the impurities in the source/drain regions 58 are activated by RTA (rapid thermal annealing) and a MOSFET with low source/drain resistance is obtained (see FIG. 2F).\nIn the above-described conventional method, when the width of a field oxidation region (i.e., the distance between adjacent devices (Wi in FIG. 2B)) is reduced to 0.2 xcexcm or less (xe2x80x9csub-quarter micronxe2x80x9d), there arises a problem of insufficiency of an oxidation amount in the dry oxidation process and a resultant insufficiency in thickness of the thermal oxidation film. One of the reasons for this insufficiency in the oxidation amount is stress generated in the SOI substrate at the time of forming the openings for the field oxidation regions (in the step b).\nTo obtain a sufficient amount of oxidation, an oxidizing temperature may be increased and oxidizing time may be lengthened. However, thermal oxidation at a high temperature for a long time will cause stress in the whole SOI substrate (i.e., in the wafer). This stress may induce defects in crystals in the substrate or cause warping of the substrate. Thus, if the oxidation is conducted at high temperature for a long time to ensure a sufficient amount of oxidation in hyperfine isolation regions of about 0.2 xcexcm, the amount of oxidation will be excessively increased at areas where the design rules are less strict (e.g., peripheral circuits); i.e., the device isolation regions at those areas may be relatively wide. The thickness of the silicon layer of the SOI substrate is thinner than the conventional silicon substrate (silicon wafer). For example, the typical thickness of the silicon layer of the SOI substrate is about several nm, while the typical thickness of the conventional silicon substrate is, for example, about 625 xcexcm. Therefore, the increase of amount of oxidation may significantly cause stress in the peripheral circuit regions of the LSI, in particular, formed in the SOI, and thus cause increases in leakage currents, for example. Such effects may adversely affect the operating characteristics of the LSI which is formed on an SOI substrate.\nIn view of the aforementioned, an object of the present invention is to obtain a sufficient amount of oxidation, without changing oxidation conditions such as temperature or time, during forming of device isolation regions of 0.2 xcexcm or less by thermal oxidation.\nTo achieve the above object, a first aspect of the present invention is a method of fabricating a MOSFET, the method comprising:\n(a) preparing an SOI substrate;\n(b) depositing an oxide film on the SOI substrate;\n(c) depositing a nitride film on the oxide film;\n(d) forming an opening in the nitride film and oxide film at a predetermined region, at which a device isolation region is to be formed, by lithography for exposing a surface of the SOI substrate;\n(e) irradiating the substantially the entire area of the silicon substrate with Ar ions;\n(f) forming a field oxide film by dry oxidation; and\n(g) removing remaining portions of the nitride film and the oxide film.\nIn a second aspect of the present invention, Si ions are used in place of the Ar ions in the first aspect.\nA third aspect of the present invention is a method for fabricating a MOSFET, the method comprising:\n(a) preparing an SOI substrate having a structure of silicon layer/buried oxide/substrate;\n(b) depositing an oxide film on the SOI substrate;\n(c) depositing a nitride film on the oxide film;\n(d) forming an opening in the nitride film and oxide film at a predetermined region, at which a device isolation region is to be formed, by lithography for exposing a surface of the SOI substrate;\n(e) irradiating substantially the entire area of the SOI substrate with at least one of Ar ions and Si ions for implanting the at least one of Ar ions and Si ions into the silicon layer of the SOI substrate in the vicinity of the surface exposed by the step of forming the opening, the nitride film and the oxide film serving as a mask;\n(f) forming a field oxide film by dry oxidation; and\n(g) removing remaining portions of the nitride film and the oxide film.\nIn each aspect, the thickness of the oxide film is preferably about 5-10 nm, and the thickness of the oxidation-resistant nitride film provided on the oxide film is preferably about 50-150 nm. The ion implantation is preferably conducted at an implantation energy of 40-50 keV, and implantation dose of 1xc3x971014 to 5xc3x971015 cmxe2x88x922.\nThrough the ion implantation under these conditions, the regions of the substrate where the openings are formed become amorphous, while defects in the substrate at the regions where devices are to be mounted can be avoided. Therefore, the field oxidation is enhanced, and the thickness of the thermal oxidation film will be sufficient even at the device isolation regions having openings of 0.2 xcexcm or less. Further, no harmful effects will be caused to the electric characteristics of the device."}
{"text": "According to the prior art of a wrench holder disclosed in U.S. Pat. No. 5,346,063, it comprises a flat main frame which has two pieces of plates protruded from its both sides on the front. There is a plurality of recesses symmetrically defined on the two plates for wrenches to be placed. Each opening of the recess has a hook, and the wrenches are prevented falling from the recesses by the hooks. There are certain disadvantages from the prior art:\n1. Wrenches are not able to be firmly held in the recesses with only two plates and simple-designed recesses on the main frame.\n2. The holder can only be laid down due to the flat design of the main frame. Because the wrench holder is not standable, it causes inconvenient for users to choose proper wrenches."}
{"text": "1. Field\nMethods and apparatuses consistent with one or more exemplary embodiments relate to a method of displaying information or a user interface (UI) by a device, and the device, and more particularly, to a method of displaying appropriate information or an appropriate UI on a user device and the user device.\n2. Description of the Related Art\nWhen using various appliances such as a mobile phone, a smartphone, a laptop computer, a tablet personal computer (PC), a handheld PC, an electronic book terminal, a digital broadcasting terminal, a personal digital assistant (PDA), a portable multimedia player (PMP), a navigation device, or a smart television (TV), a user may arrange a widget or an application execution icon on a background screen or a home screen.\nHowever, according to the related art, a user background screen or a home screen of a user device is fixed regardless of information desired by a user, thus providing unwanted information or an unwanted UI to a user."}
{"text": "Development of substances used in a variety of applications often requires an understanding of how the substances move through materials. For example, an ability of a substance (e.g., drugs, chemicals treatments, and various particulates) to diffuse through a semi-permeable material construct can provide insight into an effectiveness or a toxicity of the substance, as well as characteristics of the material construct. In some implementations, diffusion cells can be used to examine such parameters."}
{"text": "This invention concerns the temperature measurement of hot gases within a vessel having a hot surface which radiates microwave energy. The method is particularly suited to gases laden with entrained solids.\nGenerally temperatures inside of furnaces, reactors, incinerators, and the like are measured by optical or infrared pyrometers since the temperatures involved may exceed the capabilities of, e.g., thermocouples. These optical or infrared pyrometers are aimed at the point of interest through a sight hole or inspection door. In the case of modern coal gasification reactors which operate under significant pressure, transparent windows of quartz and the like, can be used. However the harsh environment within the reactor makes it extremely difficult to keep these transparent windows clear for extended periods.\nMany of the problems of such prior art are overcome by the use of a microwave radiometer to sense microwave energy emitted from a heated vessel. As described, e.g., in U.S. 4,568,199, incorporated herein by reference, the sensed energy is converted into a signal indicative thereof, and the amplitude of the signal is measured as an indication of the general temperature inside the vessel. This signal measures not only the temperature of the gases within the vessel, but is influenced by the temperature and microwave absorption properties of any suspended solids, and also by the radiation from any refractory or slag lining in the vessel which will have a temperature which will decrease as the distance increases from the hottest surface of the refractory. These influences tend to dampen the ability to measure rapid temperature fluctuations as may occur within the vessel, not only with regard to actual temperature of the gases, but also in creating a time lag as the temperatures within the entrained solids and in the refractory equilibrate with the temperature of the gases at their surface. For some applications, such as chemical conversion processes, it would be highly desirable to obtain a signal that would indicate the temperature of the gases and just the surface, of any entrained solids in contact with such gases (which surfaces would virtually be at the same temperature), and wherein such signal is capable of tracking rapid fluctuations in said temperature. Such a method has now been found."}
{"text": "The present invention is a continuation-in-part application of my application Ser. No. 750,171, filed Dec. 13, 1976 and now abandoned.\nThis invention relates generally to portable containers for food and more particularly to an insulated, pliable lunch bag adapted for carrying food and drinks while protecting them from deterioration.\nPortable containers for carrying food have been known for many years. The most popular containers still widely in use today are lunch boxes. Containers of this type, while somewhat satisfactory in some instances, are characterized by a disadvantage in that their body is of a rigid metallic or plastic construction and in that the food items placed therein are frequently susceptible to deterioration or soft drinks become undesirably warm in a relatively short period of time, especially in hot weather. Other types of food containers in form of bags or the like receptacles constructed from fabric or plastic materials which constitute the closest prior art of which I am aware have been described in U.S. Pat. Nos. 2,289,254 to Eagles and 2,667,198 to Klein. However, such bags have the common disadvantage of lacking pliability necessary for being folded or rolled up when empty. Consequently their overall size remains substantially the same after the food items have been removed which renders such bags somewhat cumbersome and inconvenient as they require to be hand-carried in empty condition. Moreover, the constitutional features of such bags are distinct from those of the lunch bag of this invention."}
{"text": "1. Field of the Invention\nThe present invention relates generally to packaging of electronic devices in the form of semiconductor dice. More particularly, the present invention relates to embodiments of an interposer for mounting a semiconductor die, wherein the interposer includes flexible solder pad elements configured for attachment to a carrier substrate or to the semiconductor die. The present invention further relates to materials and methods for forming the interposer.\n2. State of the Art\nAn electronic device in the form of a semiconductor die or chip is conventionally manufactured of materials such as silicon, germanium, or gallium arsenide. Circuitry is formed on an active surface of the semiconductor die and may include further levels of circuitry within the die itself. Due to the materials used and the intricate nature of construction, semiconductor dice are highly susceptible to physical damage or contamination from environmental conditions including, for example, moisture. In order to protect a semiconductor die from environmental conditions, it is commonly enclosed within a package that provides hermetic sealing and prevents environmental elements from physically contacting the semiconductor die.\nIn recent years, the demand for more compact electronic devices has increased, and this trend has led to the development of so called “chip-scale packages” (CSPs). One exemplary CSP design is typified by mounting a semiconductor die to a substrate, termed an interposer, having substantially the same dimensions as the semiconductor die. Bond pads of the semiconductor die are electrically connected to bond pads on a first surface of the interposer, and the semiconductor die is encased within an encapsulant material. Conductive pathways, which may comprise a combination of traces and vias, extend from the interposer bond pads to a second, opposing side of the interposer where they terminate in external electrodes to which further electrical connections are made. Typically, a CSP is then mounted to a carrier substrate, such as a circuit board having a number of other electronic devices attached thereto.\nElectrically connecting the bond pads of a semiconductor die to the bond pads of a CSP interposer generally involves using one of two types of interconnection methods, depending on the manner in which the semiconductor die is mounted. As shown by FIG. 1, a CSP 2 is configured with the first interconnection method by mounting a semiconductor die 4 to an interposer 6 with die bond pads 8 in a face-up orientation, and electrically connecting die bonds pads 8 to interposer bond pads 10 with bond wires 12. As shown by FIG. 2, CSP 2′ is configured with the second interconnection method by mounting semiconductor die 4 with die bond pads 8 in a face-down or flip-chip orientation, and electrically connecting die bond pads 8 directly to interposer bond pads 10 with conductive elements, such as bumps 14, formed of solder or a conductive adhesive material. Once the interconnection method used for CSP 2 or CSP 2′ is complete, semiconductor die 4 is encased within an encapsulant material 15 such as a polymer-based molding compound.\nFurther, FIGS. 1 and 2 show there are generally two types of external electrode structures used for mounting CSPs to a carrier substrate 16. CSP 2 of FIG. 1 is configured as a land grid array (LGA) type package, wherein the external electrodes comprise solder pads 18 that are intended to be directly attached to corresponding solder pads 20 on a carrier substrate 16. In FIG. 2, CSP 2′ is configured as a ball grid array (BGA) type package, wherein the external electrodes comprise solder ball pads 22 having solder balls 24 formed thereon, such that solder balls 24 will be attached to the solder pads 20 on carrier substrate 16.\nAlthough CSPs of the type described above have provided a compact and economical approach to packaging of semiconductor dice, they still present certain disadvantages, especially in terms of the LGA or BGA electrode structures used for mounting CSPs to a carrier substrate.\nDuring the operation of an electronic device configured as a CSP, for example, the functioning of the circuits within the semiconductor die and resistance in the circuit connections of the semiconductor die, interposer, and carrier substrate generate heat. This heating results in the expansion and contraction of all of these components as temperatures rise and fall. Because the semiconductor die, interposer, and carrier substrate are made of different materials exhibiting different coefficients of thermal expansion (CTE), they expand and contract at different rates during thermal cycling. This mismatch in thermal expansion rates places stress on the electrode structures joining the CSP interposer to the carrier substrate, and may eventually cause cracks in the electrode structures leading to the failure of electrical connections.\nThis thermal stress may be especially problematic with a CSP configured as an LGA type package as in FIG. 1, because the stress is concentrated within the relatively small thickness H of the solder pads 18 between interposer 6 and carrier substrate 16. With a CSP configured as a BGA type package as in FIG. 2, the thermal stress may be more effectively absorbed by being spread across the increased thickness H′ provided by the solder balls 24. However, because modem circuitry layouts tend to require increasing numbers of I/Os, the external electrodes on a CSP must be very densely spaced, and there are physical limits to the minimum spacing that may be attained when forming solder balls 24. The conventional process of printing and reflowing solder paste on solder ball pads 22 to form solder balls 24, for example, requires that solder ball pads 22 must be spaced at a pitch of about 0.4 mm to ensure reliable formation without bridging. Furthermore, high I/O CSPs require the use of smaller diameter solder balls that may not provide a thickness H′ sufficient to overcome thermal induced stress failures. Forming a CSP as a BGA type package also includes the additional processing required to form solder balls 24, which is undesirable in terms of mass-scale production.\nAnother problem associated with prior art package interposers is that the LGA or BGA type external electrode structures are typically formed entirely of metal or metal alloys and are, therefore, rigid. In many cases, one or both of the interposer and the carrier substrate to which it is to be mounted may have uneven surfaces or may become warped by thermal stresses during attachment of a CSP by solder reflow. When this occurs, the space between the interposer and the carrier substrate may vary, and the rigid construction of LGA or BGA type external electrodes in contact with the carrier substrate at narrower spaces may prevent contact by external electrodes at wider spaces.\nIn view of the foregoing, what is needed is an interposer for a semiconductor die package such as a CSP that is simple and inexpensive to produce and overcomes the problems associated with the prior art external electrode structures used to mount the interposer to a carrier substrate."}
{"text": "This invention relates to embossing devices, specifically to a method for placing embossed braille markings on currency notes by marking labels to affix to currency notes. Applications could include, but are not limited to, marking United States currency notes, the currency notes of other nations, and currency note substitutes such as checks, coupons, travelers checks, and money orders for the benefit of visually challenged people."}
{"text": "Light weapons, which are used against human beings, have been widely used for many years now. There are known a variety of ammunition for light weapons, such as: conventional bullets, tracing bullets that are mostly used under dark conditions, hollow point bullets that explode upon impact, hand grenades, shells, etc.\nHereinafter, when the terms “gun” or “rifle” are used, if not otherwise specifically stated, these should be understood as relating to any type of light weapon, such as a pistol, shotgun, hunting weapons, machine gun, automatic rifle, submachine gun, etc.\nWhen shooting a bullet of any type with a rifle, very accurate aim is required in order to damage the target, as the bullet generally does not include any charge. If the bullet misses the target, the target suffers no damage. Even when the bullet hits an object proximate to the target, the damage to the target is in most cases minor. The ricochets that a bullet can cause are generally small in size, and without most of the energy of the impact, the resulting damage to the target is small. Therefore, whoever wants to escape from a direct hit by such a bullet has to seek protection behind a strong and solid object. In such a case, if the bullet cannot penetrate this solid object, or event when it penetrates but loses most of its energy, the target generally escapes significant damage.\nThe art has failed to provide means for a light weapon firing a bullet to hit a target hidden behind a strong and solid object. The only solution that the art has provided to that problem is the firing of a hand grenade, or a heavier projectile containing explosive charge by means of a heavier weapon. Moreover, no solution has been provided yet by the art for the case in which the bullet passes very close to the target, but misses it.\nIt is therefore an object of the present invention to provide a bullet for a light weapon, which can cause significant damage to a target, even without directly hitting it.\nIt is another object of the present invention to provide a bullet that can damage a target hidden behind a strong and solid object.\nIt is still another object of the present invention to provide said bullet with no change to the structure of the cartridge or the firing weapon.\nOther objects and advantages of the invention will become apparent when the description proceeds."}
{"text": "Conventionally, as a fastener element used in a slide fastener, there have been known synthetic resin fastener elements, each being individually formed by performing injection molding of synthetic resin on a fastener tape, continuous fastener elements formed by forming monofilament in a coil shape or a zigzag shape, metal fastener elements formed by swaging a Y-shaped metal element material onto a fastener tape, or the like.\nThe synthetic resin fastener elements are usually formed to straddle a fastener so as to be disposed on a first surface as an outer surface of the fastener tape and a second surface as a tape back surface. The synthetic resin fastener elements include an upper half element portion arranged on the first surface side of the fastener tape and a lower half element portion arranged on the second surface side of the fastener tape.\nThe synthetic resin fastener element is mostly formed such that the upper half element portion and the lower half element portion have symmetrical shapes, but there is a case in which the upper half element portion and the lower half element portion are asymmetrically formed in different shapes, for example, in order to improve the appearance, the feel (the touch), or the like of the slide fastener.\nExamples of the fastener elements in which the upper half element portion and the lower half element portion have shapes different from each other are disclosed, for example, in Japanese Utility Model Application Publication No. 45-33956 (Patent Document 1), Japanese Patent Application Publication No. 47-37061 (Patent Document 2), Japanese Patent Application Laid-Open No. 2006-320642 (Patent document 3), and the like.\nFor example, a synthetic resin fastener element 51 described in Patent Document 1 includes an upper half element portion 53 arranged on a first surface side of a fastener tape 52 and a lower half element portion 54 arranged on a second surface side of the fastener tape 52 as illustrated in FIG. 12. The upper half element portion 53 includes a first tape-sandwiching portion 53a having a nearly rectangular shape in the front view and a triangular head 53b that extends from the first tape-sandwiching portion 53a toward the outside of the tape and is formed in a triangular shape in which a dimension in a tape length direction (hereinafter, this dimension is referred to as an element width dimension) gradually decreases toward the forefront thereof.\nThe lower half element portion 54 of the fastener element 51 includes a second tape-sandwiching portion that sandwiches the fastener tape 52 together with the first tape-sandwiching portion 53a, a neck 54a that extends from the second tape-sandwiching portion toward the outside of the tape and has a shape that is constricted in the element width direction, and a coupling head 54b that extends from the forefront of the neck 54a to be shaped like a bulge.\nParticularly, in the fastener element 51 of Patent Document 1, an upper surface of the upper half element portion 53 and a lower surface of the lower half element portion 54 are continuous flat surfaces, and the triangular head 53b of the upper half element portion 53 is set to be thinner in thickness (dimension of the vertical direction) than the coupling head 54b of the lower half element portion 54.\nAccording to Patent Document 1, since a slide fastener 50 is configured by using the above-described fastener element 51, the entire thickness of the fastener element 51 can be reduced. Further, since part of the coupling head 54b of the other coupling party can be covered with a side edge portion of the triangular head 53b, strong coupling to the extent that chain breaking is difficult to occur is obtained.\nNext, a synthetic resin fastener element 61 described in Patent Document 2 includes an upper half element portion 62 and a lower half element portion 63 as illustrated in FIGS. 13 and 14. The upper half element portion 62 of each fastener element 61 has a nearly trapezoidal shape when viewed from the front and includes a first tape-sandwiching portion 62a that sandwiches a fastener tape together with the lower half element portion 63, a neck 62b that extends from the first tape-sandwiching portion 62a toward the outside of the tape, and a coupling head 62c that extends from a forefront of the neck 62b to form a bulge shape. In Patent Document 2, the coupling head 62c of the upper half element portion 62 is formed to have a cross section of a semicircular shape.\nMeanwhile, the lower half element portion 63 of the fastener element 61 includes only a second tape-sandwiching portion 63a arranged to correspond to the first tape-sandwiching portion 62a and the neck 62b of the upper half element portion 62. Below the coupling head 62c of the upper half element portion 62, a space with nothing formed is present.\nIn the slide fastener 60 of Patent Document 2 having the above-described fastener element 61, since the coupling head 62c of the upper half element portion 62 has a semicircular cross section, the fastener tape 64 can be easily bent toward the upper half element portion 62 side at a small curvature in a state in which the left and right fastener elements 61 are coupled. According to Patent Document 2, a tape inner side edge portion of the first tape-sandwiching portion 62a in each fastener element 61 is formed to have a larger width than the coupling head 62c, and thus coupling of the left and right fastener elements 61 can be prevented from coming loose in the bending state of the fastener tape 64.\nA synthetic resin fastener element 71 described in Patent Document 3 has an upper half element portion 72 and a lower half element portion 73 as illustrated in FIG. 15. The upper half element portion 72 of each fastener element 71 includes a first tape-sandwiching portion 72a having a nearly rectangular shape when viewed from the front and a triangular head 72b that extends from the first tape-sandwiching portion 72a toward the outside of the tape and is formed in a triangular shape in which an element width dimension gradually decreases toward the forefront.\nThe lower half element portion 73 of the fastener element 71 includes a second tape-sandwiching portion 73a that sandwiches a fastener tape 74 together with the first tape-sandwiching portion 72a, a neck 73b that extends from the second tape-sandwiching portion 73a toward the outside of the tape, and a coupling head 73c that extends from a forefront of the neck 73b to form a bulge shape.\nFurther, in the fastener element 71, an upper surface of the upper half element portion 72 is formed to be curved such that a central portion in an element width direction and a central portion in an element length direction can protrude upward. Meanwhile, a lower surface of the lower half element portion 73 is formed to be a flat plane surface.\nIn a slide fastener 70 having the fastener element 71 of Patent Document 3, the upper surface of the upper half element portion 72 has the curved surface that bulges in the form of a gentle circular arc, and thus the feel or the touch of the fastener element 71 can be improved, and the appearance of the element row can be improved.\nPatent Document 1: Japanese Utility Model Application Publication No. 45-33956\nPatent Document 2: Japanese Patent Application Publication No. 47-37061\nPatent Document 3: Japanese Patent Application Laid-Open No. 2006-320642"}
{"text": "The present invention relates to the anaerobic digestion treatment of biochemical wastes. More particularly, the invention relates to a possess for removing BOD (biological oxygen demand) and nitrogen and phosphorus from waste water discharged from the process of the anaerobic digestion treatment of a biochemical waste.\nThe anaerobic digestion treatment of biochemical wastes have heretofore been adopted for treating wastes having a high BOD value, for example, excessive activated sludges, excretions and waste waters from the alcohol distillation process. The anaerobic digestion treatment has recently attracted attention in the art because of various advantages. For example, since the digestion gas generated as a by-product can be used as the energy source of the treatment apparatus, the treatment cost is low, and the digestion sludge discharged as a by-product can be used as a high quality organic fertilizer.\nOne example of such anaerobic digestion treatment is disclosed in the specification of U.S. Patent Application Ser. No. 685,901 filed on May 12, 1976 now U.S. Pat. No. 4,067,801, and this treatment process comprises the following two main steps; the primary treatment step of subjecting the starting waste water to pulverization, heating and acid addition to form an organic slurry, and the secondary treatment step of subjecting the organic slurry to anaerobic fermentation in the presence of bacteria to decompose it to digested sludge and methane gas. In general, since BOD is left at a concentration of several hundred ppm in the anaerobically treated water (the effluent separated from the digested sludge by solid-liquid separation), the treated water is diluted and subjected to the aerobic digestion treatment to remove BOD therefrom and then, the treated water is discharged into sewerage or the like. This discharged water still contains 500 to 800 ppm of nitrogen and about 10 ppm of phosphorus. Accordingly, the anaerobic digestion treatment involves a problem of occurrence of secondary pollution by promotion of eutrophication of rivers, seas and the like by nitrogen and phosphorus contained in discharged water. As means for removal of nitrogen and phosphorus, there are known physical, chemical and biological processes, but there is not known a process for removing nitrogen and phosphorus from waste waters from the anaerobic digestion treatment process."}
{"text": "Computer networks, such as the Internet, enable transmission and reception of a vast array of information. In recent years, for example, some commercial retail stores have attempted to make product information available to customers over the Internet. It is becoming increasingly popular for information providers to provide mechanisms by which customers can compare such product information across multiple manufacturers and retailers. For simplicity, manufacturers, retailers, and others that sell products to customers are interchangeably referred to herein as “merchants.” For example, Internet search/shopping sites allow customers to compare pricing information for products across multiple merchants.\nCustomers searching for a product often desire to view products that are related to or serve as alternatives to that product. For example a customer searching for a particular tablet computer may want to view product information regarding comparable tablet computers offered by another manufacturer. However, conventional methods for identifying alternative products are inadequate and often require a significant amount of time to gather enough data to predict with confidence that two products are alternatives."}
{"text": "The present invention is related to the field of silica crucibles, and more specifically to a silica crucible having a multi-layer wall in which one or more of the wall layers are doped with aluminum.\nSilicon wafers used in semiconductor industries are made from ingots of single-crystal silicon. Such ingots generally are manufactured by one of two processes: the Czochralski (CZ) process and the floating zone (FZ) process. Among those, the CZ-process is more widely used for mass production of single-crystal ingots.\nIn the CZ-process, metallic silicon is charged in a silica glass crucible housed within a susceptor located in a crystal growth chamber. The charge is then heated by a heater surrounding the susceptor to melt the charged silicon. The susceptor typically is rotated during this procedure. A single silicon crystal is pulled from the silicon melt at or near the melting temperature of silicon.\nFor higher ingot productivity, a more rapid CZ-process is desirable. However, accelerating the crystal pulling rate beyond a certain rate results in improper silicon crystal structures. Many trials are done to shorten the “melt down” period by increasing heating power. Similarly, in the production of large-diameter wafers, the amount of silicon charge and the melt-down time are increased and more intense energy is input. The total process time is much longer than that for small-diameter ingots.\nThis harsh melt-down procedure increases the rate of crucible inner surface roughening. Compressing the meltdown period also negatively affects the rigidity of silica crucible. Silica glass is not hard enough to prevent sagging of the side wall in harsher melt down processes. A more dimensionally stable crucible is desired.\nAt operating temperatures, the inner surface of a silica crucible reacts with the silicon melt. In many cases, it is the inner layer of the crucible that undergoes a change in morphology. More exactly, the inner surface roughens after prolonged operation in the CZ-process. This roughening can cause a loss of crystal structure of the pulled ingot. Inner surface roughening renders the crucible unfit for use in silicon ingot manufacture.\nAdditionally, the inner layer of a silica glass crucible can be dissolved into the silicon melt during the CZ-process. Silicon and oxygen, the main components of a silica crucible, are not deleterious to the silicon melt. However, impurities in the inner layer of the crucible can be transferred to the silicon melt during this process. To keep the silicon melt free from such impurities, a crucible is required to be extremely pure or to be insoluble by the silicon melt.\nA standard method for making a silica glass crucible is disclosed in U.S. Pat. No. 4,935,046. Quartz grain is supplied in a rotating mold in a crucible shape. The grain is then heated by an electric arc to fuse the inner part of the formed grain, leaving the outside grain unfused. During fusion, additional grain is supplied to the inside surface of the crucible. Quartz grain is melted and piled up as a transparent inner layer, while the formed grain is fused rather promptly to make an translucent silica glass substrate. The resultant crucible has a wall comprising a transparent inner layer and an translucent outer layer having a rough outer surface, which is the interface between fused grain and unfused grain.\nOne crucible is known to reduce the dissolution of the inner surface of the crucible. U.S. Pat. Nos. 5,976,247 and 5,980,629 disclose the creation of a “devitrified” layer inside of a crucible to prevent particulate generation at the silica-melt interface. The devitrified layer is reported to dissolve uniformly. Here, the devitrified layer means a crystallized silica layer, which the present inventors found to dissolve more slowly in the silicon melt than does amorphous silica. The above references claim alkaline-earth elements as a devitrification promoter, with barium recited as an example.\nOne of the present inventors filed Japanese Patent 3100836 (laid open Tokukai Hei8-2932), teaching an inner layer containing from 100-2000 ppm aluminum and 0.5-1 mm in thickness. The inner layer crystallizes in the CZ-process, so dissolution is suppressed and the dimensional stability of the crucible is improved.\nIt is known in the ingot manufacturing industry that circular patterns (“rosettes”) are observed on the crucible surface contacting the silicon melt. Examples are shown in U.S. Pat. No. 4,935,046, FIGS. 6A-6B. The ring is referred to in U.S. Pat. No. 4,935,046 as crystobalite. This phenomenon was investigated and determined to be a rosette surrounded by crystobalite.\nThe crystobalite ring is normally decorated with brown deposit when cooled down after a CZ-process use. It is hypothesized that the brown deposit is either silicon mono-oxide or colloidal silicon. The center of the rosette has a rough surface that is either not covered by crystobalite or covered by a very thin crystobalite layer. The outside of the rosette is the original silica glass surface, which has retained its original smoothness.\nAs CZ-process time continues, rosettes grow and the surfaces of the rosette centers become rough. Further, the rosettes merge and the rough surface area is increased. The smooth virgin surface decreases and finally disappears. When a major portion of the inner surface of the crucible is covered by a rough surface, the pulled silicon crystal loses its crystalline structure. Such a roughened crucible is unsuitable for ingot manufacture and silicon crystal pulling using a roughened crucible must be ceased to avoid manufacture of substandard ingots.\nA method to reduce roughening of the crucible inner surface is disclosed in U.S. Pat. No. 4,935,046. The reference further mentions that growth of crystobalite is suppressed, as the result of applying the method. By applying the layer-by-layer deposition method as taught by this reference, however, suppression of crystobalite is insufficient and roughening still proceeds to a significant extent. “Devitrification” of the outer layer of a crucible is disclosed in U.S. Pat. Nos. 5,976,247 and 5,980,629. By coating a crucible with barium-containing chemicals, the outside of the crucible is “devitrified”, i.e., crystallized, when used in a CZ-process. This crystallized layer reinforces the crucible at operating temperatures and prevents sagging of the crucible side wall.\nBy using barium as a crystallization promoter, the crystallized layer grows as CZ-process time elapses. The silica glass experiences a large volume change when it crystallizes, creating stress at each of the interfaces of the glassy phase and crystalline phase. This stress is relieved by micro-scale deformation of the crucible. If the crystalline layer thickness exceeds a certain level, the crucible is prone to cracks and possible leakage of the silicon melt. Even if the amount of barium-doped material is carefully optimized to the running conditions, the crucible nevertheless occasionally experiences cracking toward the end of a CZ-process run.\nJapanese Patent P2000-247778A discloses a three-layer crucible. The layers are a transparent synthetic silica inner layer, a synthetic silica or natural quartz glass middle layer, and an aluminum-doped silica outer layer. The optimum range for aluminum concentration in the outer layer is reported to be 50-120 ppm. The best mode taught in this reference has an approximately 3 mm thick outer layer doped with aluminum at about 75 ppm.\nThe doped aluminum outer layer helps to prevent sagging of crucible. However, the inner layer of this crucible is still prone to uncontrolled rosette formation and growth during a CZ-process.\nA long-life crucible is therefore desirable, especially a large-diameter CZ-process crucible. Specifically, the side wall of the crucible should be able to maintain its structural integrity without warping, and the inner surface of the wall should resist roughening during a CZ-process."}
{"text": "Many modern electronic components are created by thin film wafer processing. One category of component created by thin film processing is the tape head. Another category is the disk head.\nMost tape heads are currently built on wafers using thin film processes, similar to the wafers used for fabricating disk heads. However, the operating efficiency of disk heads and tape heads are inherently different. Disk recording/reading is very efficient, as the disk media is extremely flat and smooth, has a very thin magnetic layer, is in a sealed environment, and the heads are constructed to function with a particular media. Writing and reading tapes must address very different challenges. For example, the head must work with different tape brands, which can have different physical and magnetic properties. Furthermore, most tape is composed of magnetic particles, which are coated onto the tape surface. The resulting media can have variations in coating thickness and particle dispersion. This, coupled with spacing loss variations due to embedded wear particles and debris, requires that magnetic bits in tape must be much larger than bits in disk media for achieving an acceptable signal-to-noise ratio. Tape bits are typically of the order of 100 times wider and 3 times longer than bits recorded onto disks. Disk drive heads are designed to fly over smooth disk surfaces in a controlled manner at speeds approaching 30 to 40 meters per second. By contrast, tape stacking and other requirements limit tape drive operating speeds to approximately 3 to 6 meters per second. Thus, to achieve data rates commensurate with disk drives, high performance linear tape drives typically employ heads having multiple pairs of write-read heads that operate simultaneously. For example, two pairs provide twice the data rate of one pair, and Linear Tape Open (LTO) heads have eight pairs of read and write elements for each direction.\nOften in tape head fabrication, head images are laid out on the wafer such that the heads cut from the wafer are the required length for insertion in a tape drive. However, for LTO heads for example, the active area of the head is approximately 7 mms long, whereas the tape supporting surface of the head must be 23 mm long. The remaining 16 mm are blank, i.e. devoid of devices.\nA problem is that thin film wafers are of a standardized size, and thus, the number of individual dies which contain the read/write recording devices that can be cut from each wafer is limited. Increasing use of more complex wafer processes, coupled with the high cost of wafer substrates makes achieving the highest number of heads possible from a single wafer an important head design priority.\nAn approach to increasing the number of dies per wafer is to make each die no larger than the active area needed for each head, Thus, for example, three partial span images can fit in the space of one full span LTO image. The dies are fabricated into chips which are then inserted into a passive carrier constructed of similar substrate material. In this way, the tape is fully supported over the width of the head, but wafer costs are dramatically reduced. Partial span heads are conventionally fabricated such that the closure portion of the head is completely surrounded by the carrier. The entire structure, including head and carrier, are machined together to form a uniform tape bearing surface, which is planer or cylindrical and has no steps, discontinuities or corners. Thus, the carrier fully encompasses the original chip and the seam between the chip and carrier are not discontinuous. This approach generally mandates that the chip image be tall enough for mechanical processing as well as provide enough material for forming the contour. This takes away wafer space, resulting in fewer chips per wafer, etc.\nA second method is to fabricate the partial span chip, i.e. lap it, etc., so that its tape bearing surface is completed prior to assembly. This otherwise finished chip is then attached to a beam such that the tape bearing surfaces of both chip and beam has minimal discontinuities. Also, the closure portion of the chip is aligned with an inner surface of the beam to eliminate any discontinuities along the edge of the tape bearing surface. The object of this method is to make the chip and beam form as closely as possible a single, regular surface, even though in general there will be steps at the chip edges. This approach allows short partial span images, so wafer utilization is good. However, a problem inherent in this method is that it is difficult to assemble the chip and beam with near-perfect alignment at the tape bearing surface. It is also more difficult to attach a cable since the contacts are now recessed from the side of the beam, and conventional cable bonding tooling does not easily reach into the recess. It would be desirable to create a partial span head in order to obtain a high utilization of useful circuitry per wafer, thereby minimizing fabrication costs and decreasing the cost per unit of magnetic heads. It would also be desirable to avoid having to provide the head contour after the chip is cut from the wafer, as required by current partial span heads, as this not only forces the use of a taller image for handling purposes, it also makes stripe and throat height control more difficult."}
{"text": "Most electronic equipment, and in particular computers, utilize a series of chips which are connected to a motherboard in order to form the signal processing part of the equipment. Various chips may assume a single function or multiple functions which are used by the equipment. The group of chips used together is sometimes referred to as a chipset.\nFIG. 1 is a block diagram showing the arrangement of a chipset on a motherboard for a computer. The chip set 100 includes a first chip 102 which carries the central processing unit for the device. Memory controller hub 104 acts as a central controller to move data into and out of memory and to other related chips. Chip 106 is a graphics chip which generates various graphic arrangements for display. Chip 108 is the memory itself, either RAM or ROM memory. Chip 110 is an input/output controller hub which transfers data to various input/output devices. Chip 112 includes connections to a hard disk drive. Chip 114 is a chip which connects to other peripheral components.\nTypically, each chip in a chip set is formed of two parts. The first part is the core which is the circuitry which handles the main function of the device itself. Also on the chip are input/output circuits for connecting the core to other chips. For example, the memory controller 104 would have a central core and an input/output device connected to each of the four other chips 102, 106, 108 and 110 to which it is connected.\nFor every pair of chips that are connected, an interface is provided to connect the input/output devices of the chips to each other. Thus, the CPU 102 and memory controller hub 104 are connected by a front side bus (FSB) 116. Likewise, memory controller hub 104 is connected to graphics chip 106 through the advanced graphics port (AGP) 118. Memory 108 is connected to the memory controller hub 104 by a system memory bus 120. Memory controller hub 104 is connected to the input/output controller hub 110 through hub link 122. The input/output controller hub 110 is connected to the hard disk drive 112 through IDE 124. The I/O controller hub 110 is connected to the peripheral components chip 114 through the peripheral components interface 126.\nFIG. 1 also shows a clock circuit 113 which is another chip connected on the motherboard. This clock provides clock signals of various frequencies to the various other chips. These particular connections are not specifically shown but all chips on the motherboard are connected thereto to receive clock signals which are necessary for the synchronization of the entire device.\nSome of the interfaces on the motherboard are considered to be source synchronous interfaces. In the present example, the front side bus 116, the advanced graphics port 118 and the hub link 122 are all source synchronous circuits. On the other hand, a system memory bus 120 and IDE 124 are not source synchronous interfaces. In such an interface, data signals and strobe signals are used to transfer data in a synchronous fashion. These signals occur in a certain preset timing relationship so that data being transferred can be expected at a particular time location."}
{"text": "Shuffling machines, or shufflers, are widely used in casinos, card rooms and many other venues at which card games are played. Conventional shufflers are typically adapted to receive one or more decks of standard playing cards to be shuffled. The intended purpose of most shufflers is to shuffle the playing cards into what is believed to be a random order. Such a random order of the playing cards is desirable when playing various types of card games such as blackjack, poker and the like. However, in reality most shufflers have tendencies to shuffle or reorder the deck or decks in a manner which skilled card counters can perceive and use to their advantage versus the casino, house or other player. Thus, there is still a need for automated shufflers that function in a manner which more truly randomizes the ordering of a deck or decks of playing cards.\nOther problems associated with at least some conventional shufflers include excessive size, excessive weight, excessive mechanical complexity and/or electronic complexity. These complexities also may fail to achieve a suitable degree of shuffling, reordering or recompiling into a truly random order from one shuffling process to another. Accordingly, there is still a need for improved automated shuffling machines for playing cards that produce reordering of card decks in a manner which is closer to true randomness and which is more difficult for skilled card players to decipher to change the odds so as to be relatively favorable to the player versus unfavorable portions of a deck or decks of cards.\nOne casino game commonly called “blackjack” or “21” is known to be susceptible to card counting and casinos are routinely spending significant amounts of money trying to prevent card counters from taking advantage of non-random sequences in the decks held within a dealing shoe that holds the decks being dealt. Poker has also grown in popularity and is played with a single deck, which makes any knowledge of cards of potential significance to a player.\nThe inventions shown and described herein may be used to address one or more of such problems or other problems not set out herein and/or which are only understood or appreciated at a later time. The future may also bring to light currently unknown or unrecognized benefits that may be appreciated, or more fully appreciated, in association with the inventions shown and described herein. The desires and expected benefits explained herein are not admissions that others have recognized such prior needs, since invention and discovery are both inventive under the law and may relate to the inventions described herein."}
{"text": "The availability of 7 GHz of unlicensed spectrum in the 60 GHz band offers the potential for multi-Gigabit indoor wireless personal area networking (WPAN). Applications that require large bandwidth include uncompressed High Definition (HD) video streaming, fast file download from an airport kiosk (Sync & Go) and wireless display and docking, to name just a few. These applications cannot be supported over existing home networking solutions (IEEE 802.11a/b/g/n and WiMedia UWB) because the required data rates far exceed the capabilities of these networks.\nThus, a strong need exists for improvements and new development in wireless personal area networks that operate in the 60 GHz band.\nIt will be appreciated that for simplicity and clarity of illustration, elements illustrated in the figures have not necessarily been drawn to scale. For example, the dimensions of some of the elements are exaggerated relative to other elements for clarity. Further, where considered appropriate, reference numerals have been repeated among the figures to indicate corresponding or analogous elements."}
{"text": "A virtual machine system sometimes consists of multiple physical machines and runs multiple hypervisors on a single machine. Each hypervisor can support multiple virtual machines, with each virtual machine running a guest to perform tasks for a user. From time to time a system administrator may want to move (“migrate”) a guest from one hypervisor to another for maintenance or performance reasons. The migration may be a “live migration,” which means that the guest can be moved without disconnecting its client or application.\nWhen a guest migrates to a different hypervisor, its network location is considered as changed. A changed network location means that the different hypervisor (“target hypervisor”) is now responsible for forwarding packets to the guest. Switching components (also referred to as “network devices”) in the network to which the target hypervisor is coupled are notified of the guest's new association with the target hypervisor so that the guest can continue to receive packets after migration.\nOne current approach is to have the source hypervisor send the guest's network addresses to the target hypervisor. In response, the target hypervisor sends one or more notification packets to the network devices, notifying the network devices of the guest's network addresses. Thus, when a network device receives a packet destined for any of these addresses, the network device can forward the packet to the target hypervisor, which then relays the packet to the guest. Conventionally, each notification packet sent from the target hypervisor to the network devices contains one network address of the guest. Since a guest can have multiple different network addresses (e.g., in the tens or hundreds), such notification can take tens or hundreds of packets to complete. These notification packets can burden the network and reduce available bandwidth for normal network traffic."}
{"text": "1. Field of the Invention\nThis invention relates generally to vehicle windshield wiper systems and, more particularly, to a system wherein the radial reach of the wiper blade support arm varies during the angular stroke of the arm.\n2. Description of the Prior Art\nEfforts to simplify vehicle windshield wiper systems have led to proposals wherein a single wiper blade support arm is centrally supported below the windshield and oscillated through a wiping arc while the radial reach of the arm is varied to wipe toward the outer corners of the windshield. In one such proposal a wiper arm is slidably mounted on an angularly oscillating carrier and radially projected by a gear driven crank arm. In another proposal, a wiper arm is mounted on an angularly oscillating carrier by a parallelogram linkage and is radially projected by a gear driven crank arm and connecting link arrangement. In still another proposal, a radially shiftable wiper arm on an oscillating carrier is projected between its radial positions by a planetary gear driven crank arm and connecting link arrangement. A vehicle windshield wiper system according to this invention represents an economical alternative to these and other known single arm, variable reach systems."}
{"text": "Prior art methods for nucleic acid extraction from cellular source materials and, particularly, paraffin embedded tissue samples (e.g., formalin-fixed paraffin-embedded samples: FFPE) involves complicated, multi-step processes.\nThe extraction of nucleic acids from Mycobacteria in sputum, for example, is a challenge because sputum is very viscous and not easily processed for nucleic acid extraction. Sputum samples are typical solubilized using N-acetyl-L-cysteine-sodium hydroxide (NALC-NaOH) treatment (Coulter and Charache, Sputum digestion/decontamination for Mycobacteriology culture—Guidelines, SMILE, John Hopkins University, 2008) and the mycobacteria are pelleted by centrifugation. NALC-NaOH treatment does not kill the Mycobacteria and further treatment by heat and/or chemicals is done to inactivate the samples. Nucleic acids can be extracted from the cell pellet using several techniques to lyse cells. Sonication (Colin, et al., Method and apparatus for ultrasonic lysis of biological cells, U.S. Pat. No. 6,686,195, 2004), bead beating (Melendes, et al., Cell disrupting apparatus, U.S. Pat. No. 5,464,773, 1995), enzymes (Salazar and Asenjo, Enzymatic lysis of microbial cells, Biotechnol Lett (2007) 29:985-994), mixing (vortexing), mechanical shearing and chaotropic solutions (Das, et al., Method for detecting pathologenic mycobacteria in clinical specimens, U.S. Pat. No. 7,638,309, 2009) are some of the methods used to break open the pelleted cells for nucleic acid extraction. These steps are in addition to the actual extraction procedures and add complexity to and time to the entire process.\nThe extraction of nucleic acid from yeast is also one of the more challenging techniques in nucleic acid (e.g., DNA) sample preparation. Yeast are fungi and have cell walls that are difficult to lyse (Lipke and Ovalle, Cell wall architecture in yeast: new structure and new challenges, J Bacteriol 1998, 180(15):3735). Lysis buffers using chaotropic salts and detergents or alkali lysis protocols of the prior art are not very effective in lysing yeast cells directly but are used with additional steps. These additional steps can be divided into two main groups: physical methods and enzymatic methods. The physical methods can include sonication of cells (U.S. Pat. No. 6,686,195) with or without the presence of grinding particles, high powered agitation with grinding particles (U.S. Pat. No. 5,464,773) (bead beating, ball mills) or the use of high pressure mechanical shearing (e.g., French pressure cell press, as is known in the art). Enzymatic methods rely on particular enzymes such as zymolase (Salazar and Asenjo, ibid; U.S. Pat. No. 5,688,644) to weaken the cell walls such that the cells can be lysed by more conventional techniques.\nThe extraction, enrichment and isolation of nucleic acids from FFPE material is a very complicated process that requires the deparaffinization of the tissue with organic solvents, the digestion of the tissue with protease and then the extraction of the nucleic acids from the tissue. These prior art processes use multiple solutions and multiple steps. The organic solvents used are not usually miscible with aqueous solutions.\nThus, what is needed are compositions and methods that permit the efficient extraction, enrichment, isolation and purification of nucleic acids from cellular source materials, particularly Mycobacteria, yeast and FFPE samples."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a control apparatus for a braking system. More specifically, the present invention relates a braking force control apparatus that generates a braking force by pushing and separating a brake pad toward and apart from a brake rotor through a driving device.\n2. Background Information\nIn a control apparatus for a braking system, a driving state of a motor is controlled by detecting a depressing force exerted on a brake pedal by a driver, such that the braking force that corresponds to the depressing force is generated. One of the methods of estimating the braking force is disclosed in the SAE 1990-01-0482. This braking force estimation method utilizes the rigidity of the caliper that supports pads. The braking force is estimated based on the amount of displacement of the motor from a point where the rotor and the pads initiated a contact. More specifically, since there is a certain relationship between the amount of displacement of the motor and the braking force determined by the caliper rigidity once the pads and the rotor start contacting, the braking force can be controlled by controlling the amount of motor displacement. However, the amount of motor displacement includes a portion that eliminates the clearance between the pads and the rotor. In other words, the amount of motor displacement includes the portion in which the motor moves in the pressuring direction until the pads and the rotor contact each other, and a portion in which the motor moves after the pads and the rotor start contacting. The aforementioned relationship between the braking force and the amount of motor displacement determined by the caliper rigidity exists only in the latter portion. Therefore, in order to obtain from the total amount of motor displacement, the amount of motor displacement occurring after the pads start contacting the rotor, it is necessary to eliminate precisely the portion in which the motor starts moving in the pushing direction until the pads start contacting the rotor. Therefore, it is necessary to detect the precise point at which the pads and the rotor start contacting.\nOne of the control apparatuses that control the braking force by detecting the position at which the pads and the rotor start contacting is disclosed in Japanese Laid-Open Patent Application H9-137841. In this control apparatus, an axial force sensor is attached to a piston that pushes and separates the pads toward and apart from each other. While the piston is moving in a direction that separates the pads apart from each other (i.e., the braking force is released), the initial contacting point of the pads and the rotor is determined as the point at which the axial force detected by the axial force sensor is zero.\nIn view of the above, there exists a need for a control apparatus which overcomes the above mentioned problems in the prior art. This invention addresses this need in the prior art as well as other needs, which will become apparent to those skilled in the art from this disclosure.\nIt has been discovered that the preciseness of the axial force sensor disclosed in Japanese Laid-Open Patent Application H9-137841 would tends to be compromised under actual operating conditions. In particular, the preciseness of the detection of the axial force by the axial force sensor tends to be compromised because a change in property due to temperature changes and noise or degradation from vibration.\nThe present invention has been conceived in view of the aforementioned problems. More specifically, an object of the invention is to provide a control apparatus that can control braking force accurately by accurately detecting an initial contacting point of the pads (the brake friction body) and the rotor (the brake rotational body) without using an axial force sensor.\nIn order to achieve the aforementioned object, a braking force control apparatus is provided that includes a driving device, an operational status detector, an initial contact position detector and a brake control unit. The driving device is operatively configured to move a brake pad towards and away from a brake rotor. The operational status detector is arranged to detect an operational status of the driving device. The initial contact position detector is arranged to detect an initial contacting point between the brake pad and the brake rotor based on the operational status of the driving device detected by the operational status detector, when the driving device moves the brake pad toward the brake rotor. The brake control unit is operatively coupled to the driving device and the operational status detector to control the driving device. The brake control unit includes a brake command value setting portion that is configured to set a brake command value to control the driving device based on a pedal operation amount. The brake control unit further includes an adjustment portion configured to subsequently adjust the brake command value based on the initial contacting point between the brake pad and the brake rotor detected by the initial contact position detector.\nThese and other objects, features, aspects and advantages of the present invention will become apparent to those skilled in the art from the following detailed description, which, taken in conjunction with the annexed drawings, discloses a preferred embodiment of the present invention."}
{"text": "A) Field of the Invention\nThe present invention relates to a shock absorbing apparatus and method to be used in conjunction with a pile driving and/or pile pulling vibratory machine, and more particularly to such an apparatus and method which can be used effectively to isolate shocks under greatly varying load conditions imparted to the shock absorbing apparatus.\nb) Background of the Invention\nIn the construction industry, it is sometimes necessary to drive piles into the earth to provide a proper foundation for a building or other structure. One method of accomplishing this is to place the pile in a vertical position above the earth's surface and strike the upper end of the pile repeatedly with a hammer (i.e., a metal mass which is raised and dropped on the pile) until the pile has penetrated into the ground surface a sufficient distance to provide adequate bearing. A later development was to drive piles into the ground by use of a vibrating machine which oscillates the pile from zero to 20,000 cycles per minute depending on the type of machine to cause what appears to be an almost continuous motion of the pile into the earth. Under some circumstances, such a vibratory machine can cause the pile to move into the earth relatively rapidly (e.g., as fast as ten feet per second).\nA typical arrangement for such a vibratory machine is to provide a pair of weights which are mounted eccentrically for rotation about parallel axes, with the directions of rotation being opposite to one another so that the lateral forces are being cancelled out, and a net up and down vibrating force is developed by the machine. One part of the machine is coupled to the upper end of the pile, while a second part of the machine is connected through a shock absorbing device to a support member, such as a cable.\nWhen the pile is being driven into the ground, the vibratory machine is able, in large part, to act substantially independently, in that only minimal exterior support is required, this being mainly to keep the vibratory machine properly positioned. Sometimes weights are added to the shock absorbing device to provide a greater downward force, and this gives need for effective shock absorption. Another mode of operation is when a previously driven pile is being extracted from the earth, and it is necessary to impart a tension force on the pile so as to pull it upwardly. In these circumstances, a tension force (e.g., a pulling force exerted by a connecting cable) is applied through the shock absorbing device to the vibratory machine, which in turn pulls upwardly on the pile to which it is connected. The tension force exerted by the cable can vary greatly, and can vary between two tons to one hundred tons.\nFor various reasons, it is desirable that the cable be subjected to a more constant load, with the rapid vibratory loads being isolated from the cable as much as possible. However, properly isolating these vibratory loads is complicated by the fact that the tension loads necessary to extract the pile can vary greatly, depending upon the size of the pile, the depth to which it is driven, and the localized resisting forces imparted by various portions of the earth material."}
{"text": "In mobile communication and mobile computing systems, the combined need for digital and non-digital functionalities in an integrated system is driving a dual trend in system level integration.\nOne trend in system level integration is the concept of a system on chip (SOC). In a system on chip, functions (e.g., digital functions) are integrated on-chip. In these systems, the number of digital functions that may be integrated on-chip continues to increase as the number of transistors available on-chip increases. Another trend in system level integration is the functional diversification of semiconductor-based devices. Instead of implementing devices on one chip, the devices are implemented on multiple chips and integrated into one package, referred to as a system in package (SiP).\nIn a system on chip as well as a system in package, the embedded capacitors for providing the digital functions are built with metal layers. Furthermore, in a typical system, the conventional capacitors are not proximate to the active devices. This leads to larger area consumption and a more parasitic relationship. Conventional capacitor process technologies are also associated with a two-dimensional (2D) process. A FinFET process technology is associated with a three-dimensional (3D) process."}
{"text": "A variety of toothbrush head configurations exist that have manually and/or mechanically-driven movable cleaning elements. Many of these configurations, however, include cleaning elements that extend from a rigid head. Teeth and gums by nature have a complex intricate contour. Due to the rigid nature of the attachment of the cleaning elements to the head of the toothbrush, the orientation of the cleaning elements is not flexible. Thus, a need exists for a toothbrush that achieves better flexibility of cleaning elements for an enhanced and improved cleaning action during brushing."}
{"text": "Automatic darkening filters are often provided on protective headgear (e.g., headwear or eyewear), where protection from high intensity light is desired."}
{"text": "This section is intended to introduce the reader to various aspects of art, which may be related to various aspects of the present invention that are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present invention. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nDigital data (e.g., photos, videos) are increasingly produced and managed on mobiles devices (e.g., smartphones, tablets, laptops). This data is also often shared, backed up, or processed via Internet. Indeed, a wide range of “cloud” services handle users' content, be they photo processing services, social networks or online storage. Most of these cloud services rely entirely on web technologies. As a consequence, users need to upload large amounts of content over HTTP to web applications. However, the speed of uploads is limited by the available bandwidth. Indeed, the connectivity speed to the\nInternet remains limited due to the use of legacy infrastructures (xDSL), or of shared medium (Cellular).\nThe long upload times prevent users from standing by or powering off their stand-alone devices and require these users to keep their devices connected to handle the transfer over the Internet. In order to alleviate these issues, a mechanism to offload uploads over HTTP to a third party device which is permanently connected to the network, such as residential gateway, is proposed. A method for locating the third party device offering the offloading web service is therefore proposed.\nHowever, offloading a task to a third party requires to trust this third party, in other words the third party device hosting the offload service has to be authenticated by the user stand/alone device. Known solution for authenticating a device or a web service are based on certification by a trust authority. Certificates are either delivered by a trust authority to a trusted operator owning the web service or to user's physical device. However, these solutions are not compatible with the legacy software, such as web browser, and standard web protocols wherein the processing environment is limited. In others words, the browser is limited in term of inputs and outputs, for instance the browser cannot access to the storage media (such as hard disk drive) of the device on which it is executed, cannot access directly to the network.\nA solution for securely accessing a web service by a browser running a web application on a user device through a network, wherein the web service is hosted by a local device is therefore required. The method should deliberately be simple to ease implementation and use, and compatible with legacy software, adapted to be implemented in JavaScript and to run in the browser.\nThe present invention provides such a solution."}
{"text": "The existing combined fireworks are ground firework products formed by combining a plurality of single cylinders and producing the effects such as sound, light, color and floating materials. The single-cylinder product generally is formed by outer cylinder, propellant powder, inner cylinder or effective powder, ignition fuse, spreading fuse, clay plug, base, etc. In manufacturing, the outer cylinders of a plurality of single-cylinder products are firstly combined and molded through the necessary processes such as cylinder rolling, mudding, bonding and ranking, fuse connecting, drying, combining, which last for one week and result in low efficiency, slow speed, and the cylinders are easily distorted and bent with nonstandard size, thus, the scale production and standardization of the firework industry is severely restricted. Moreover, the side face of the cylinder body should be manually drilled for forming the fuse hole in the procedure of fuse connecting, so that it is not easy to control the size of the hole and the distance between holes, which affects the launching time and launching effect of the product; in addition, the fuse of the formed product is at the sidewall of the cylinder body, which will easily lead to the fire hazard when firing due to the exposure of the fuse sparks."}
{"text": "The invention relates to a longitudinal seat adjustment device.\nLongitudinal seat adjustments generally have a lower rail fixed to a chassis, inside which lower rail and upper rail, to which the seat is attached, can be propelled by motor. A spindle fixed by its respective ends to the lower rail sits inside the upper rail, on which spindle a transmission fixedly coupled to the upper rail is arranged in axially movable fashion. The seat, which is located on the two upper rails arranged parallel to one another, can be propelled via a motor device that sits between the rails.\nExamples of devices for the purpose of longitudinal seat adjustments have been disclosed in DE 36 40 197 A1, DE 42 08 948 C2, DE 196 42 655 C2, DE 198 15 283 A1, DE 198 44 817 A1, DE 199 44 690 A1, and WO 95/16 585. There is a need for a compact longitudinal seat adjusting device in which only a few components are needed and that can without difficulty accommodate the forces that occur especially in the case of a crash, and in particular for a transmission that transforms the rotational motion of the drive motor into a translational motion.\nA further goal of the present invention is for the space occupied by the longitudinal seat adjusting device to be relatively small, preferably only 15 mm in width. Moreover, the permissible protrusion of the upper rail is likewise not to be too great, at most some 15 mm. Finally, the device according to the invention is to satisfy a strength requirement in both directions, which is for example 25000 N.\nFinally, the device according to the invention is to be relatively rapidly adjustable, that is, for example at between 15 and 25 mm/s over a relatively long adjustment range of for example 300 mm."}
{"text": "It is conventional, in providing aspirators in an analyzer, to sense the position of the aspirator to prevent vertical overtravel. For example, analyzers provided by Eastman Kodak Company under the trademark \"Ektachem 700\" have a sensor to determine when the aspirator has lowered a sufficient distance to pick up, by force-fit, a disposable tip. Such sensors are vertically adjusted so that, as part of the set-up calibration, the operator can adjust the trigger location of the sensor depending on that particular analyzer's margins of error. Other instruments have similarly adjusted sensors. See, for example, U.S. Pat. No. 4,705,667 (Column 6, lines 37-44). Once adjusted, the sensor then cooperates with a single flag for each movable part, since that flag is sufficient to trigger that the limiting condition is achieved.\nOne difficulty of such a construction is that when replacement parts are incorporated into the analyzer proboscis, the set-up calibration (including the positioning of the sensor) has to be recalibrated, which has to be done manually. This is time-consuming and expensive.\nTherefore, prior to this invention, there has been a need to provide a proboscis with a sensor for set-up calibration that allows recalibration without readjusting the sensor. This has not been possible with the conventional configuration described above.\nYet another drawback of the previous construction has been that, with only a single flag acting with the sensor, the proboscis can detect only that it has been interrupted. It is unable to determine the nature of the interfering objects since, although the proboscis may be stalled, the flag is not in position to interact with the sensor."}
{"text": "1. Field of the Invention\nThis invention generally relates to an image sensing device for converting light image information into electrical signals, and, in particular, to such an image sensing device which may be advantageously used as an image sensor of a facsimile machine and the like for reading the document information to be transmitted to a remote place. The present invention also relates to photoelectric elements which are particularly suited for use in an image sensing device and the method of manufacturing such photoelectric elements.\n2. Description of the Prior Art\nAs is well known in the art, an image sensing device usually includes a plurality of light-sensitive or photoelectric elements arranged in the form of a single array, and the plurality of photoelectric elements are activated in timed sequence from one end of the array to the other in repetition thereby scanning a document to be read along scanning line sectors in a stepwise manner. In such an image sensing device, it is common practice to divide the plurality of photoelectric elements into a predetermined number of blocks and to provide a common electrode for each block thereby having one end of each of the photoelectric elements belonging to the same block connected to the corresponding common electrode mainly for simplification of wiring. With such a structure, it is true that wiring mav be simplified, but each of the photoelectric elements looses an operative independency so that the image sensing device would suffer from stray signals unless an adequate measure is taken\nSuch being the case, when it is so structured to group a plurality of light-sensitive elements into blocks, each of which includes a predetermined number of light-sensitive elements, then it is necessary to apply means for preventing the occurrence of stray signals. FIG. 1 illustrates one prior art approach for preventing the occurrence of strav signals in an image sensing device, which uses a plurality of blocking diodes one for each photoelectric element As illustrated, a plurality of photoelectric elements P.sub.ij (i=1-n and j=1-m) are disposed in the form of a single array, and these elements are grouped into n blocks B.sub.1, B.sub.2, . . . , B.sub.n, each comprising m photoelectric elements. In each of the blocks, the top ends of the photoelectric elements are commonly connected to the corresponding common block terminal. For example, in the leftmost block, the top ends of the photoelectric elements P.sub.11, P.sub.12, . . . , P.sub.1m are commonly connected to the corresponding block terminal B.sub.1.\nOn the other hand, the bottom end of each of the photoelectric elements defines an individual electrode which is connected to the anode of the corresponding blocking diode D.sub.ij. Each of the blocking diodes in one block, for example the diode D.sub.11 has its cathode interconnected to the cathodes of the corresponding blocking diodes, D.sub.21, . . . , D.sub.n1 for D.sub.11, in the other blocks through respective interconnection lines. These interconnection lines are connected to ground through respective load resistors R.sub.L and also to an output terminal V.sub.0 through analog switch S.sub.1, S.sub.2, . . . , S.sub.m, respectively. Provision of blocking diodes in this manner allows to establish operative independency for each of the photoelectric elements; however, it still suffers from various disadvantages since it is necessary to provide a relatively large number (m.times.n) of diodes corresponding in number to the photoelectric elements and the manufacture of such large number of diodes at a time tends to be expensive even if use is made of the thin film forming technology and moreover it tends to lower the yield.\nFIG. 2 shows another prior art approach in which operational amplifiers are provided and their virtual ground is utilized to establish operative independency between the photoelectric elements. Similarly with the structure of FIG. 1, the top ends of the elements in each block in the arrangement of FIG. 2 are commonly connected, and the individual electrodes defined by the bottom ends of the photoelectric elements in one block are interconnected to the corresponding individual electrodes in the other blocks through the respective interconnections which are also connected to the inverting input of the respective operational amplifiers A.sub.1, A.sub.2, . . . , A.sub.m. These op amps have their outputs connected to the output terminal V.sub.0 through respective analog switches S.sub.1, S.sub.2, . . . , S.sub.m and their non-inverting inputs connected to ground. A resistor R.sub.f is connected between the output and the inverting input in each of the op amps thereby defining a feed-back loop. In such a structure, the inverting inputs of the op amps are, in effect, at virtual ground, so that the elements P.sub.ij may be operated independently one from another. However, the structure shown in FIG. 2 requires the provision of so many op amps corresponding to the number of the individual electrodes (m in the illustrated example), which also tends to push up the manufacturing cost partly because of difficulty in mounting of so many op amps.\nA further prior art approach is illustrated in FIG. 3, in which case, use is made of a mxn bit shift register for directly driving an array of photoelectric elements. In this case, however, using a mxn bit shift register requires mxn number of connections to be made, which also pushes up the cost. Such a disadvantage can not be obviated even if the shift register is constructed in the form of an IC."}
{"text": "1. Field of the Invention\nThe present invention relates to integrated circuits and, more particularly, to microwave monolithic integrated circuits (MMICs) for use in high-frequency applications, such as those in the microwave and millimeter-wave frequency regions.\n2. Description of Related Art\nIntegrated circuits have gained widespread use in many electronic applications. In early hybrid integrated circuits, active elements (such as diodes and transistors) and passive elements (such as resistors, capacitors, and inductors) were typically discrete components mounted (e.g., soldered or bonded) to a dielectric slab or substrate. In contrast, in a monolithic integrated circuit (or “monolithic circuit”), circuit components including active and passive elements are integrated monolithically, i.e., formed directly on a common semiconductor substrate.\nTypically, depending on the operating frequency, monolithic integrated circuits may be formed on different types of substrates. As an example, the monolithic integrated circuits operating up to 1-2 GHz may be fabricated on silicon (Si) substrate. At higher operating frequencies, such as microwave and millimeter-wave frequencies (approximately between 1-300 GHz), the substrate is usually gallium arsenide (GaAs) and these circuits are commonly referred to as monolithic microwave integrated circuits, or MMICs. Some of the advantages of MMICs include their small sizes, the inclusion of multiple functions (e.g., radio frequency (RF) and logic) on a single semiconductor chip, and a wider frequency-bandwidth performance that is often difficult to achieve with discrete devices due to bandwidth-limiting parasitics associated with discrete-device packaging.\nTypically, RF signals at microwave and millimeter-wave frequencies can easily penetrate harsh environments such as dust, smoke, and snow, and are very attractive due to their high spatial resolution, resulting in a compact chip size and small antenna dimension. As such, MMICs find use in various commercial, military, and space applications. For example, in addition to the traditional use in radars, microwave and millimeter wave techniques are finding applications in such diverse areas as forward-looking automotive radar, Synthetic Vision Systems (SVS) for aircraft landing, Concealed Weapon Detection (CWD) systems, industrial sensors and accelerometers.\nTypically these systems employ a stable transmitter and highly sensitive receiver incorporating a mixer and a local oscillator. However, increasing use of microwave and millimeter-wave frequency bands for communication, radar and measurements have created the need for more sophisticated methods for controlling the frequency, power and phase of these sources of radiation. For example, coherent radar systems have for years relied on phase-locked or injection-locked transmitters as well as phase-locked local oscillator for receivers. In addition to down-converters, IF amplifiers, and phase detectors, typically these techniques use a directional coupler and a power divider to meet the required specifications.\nIn general, the coupler and power divider are either coaxial/waveguide or fabricated using hybrid microstrip technology. The latter does generally reduce the overall component size, but it still does not lead to a compact, low cost design. Moreover, many newer applications require RF power monitoring capability for accurate control of output power.\nThus, there is a general need for a MMIC, which incorporates the functions of power distribution and power monitoring over a large bandwidth in the microwave and millimeter-wave frequency range."}
{"text": "1. Field\nThis application relates to an interface between a processor memory controller and a non-volatile memory controller and more particular relates to a direct interface between a memory controller and a non-volatile memory controller using a command protocol.\n2. Description of the Related Art\nIn typical computing devices, main memory includes volatile memory such as dynamic random access memory (“DRAM”) and static random access memory (“SRAM”). A processor typically communicates with the main memory over a wire interface using a low-level wire protocol such as the Joint Electron Devices Engineering Council (“JEDEC”) protocol, the industry standard for processor—DRAM interfaces. The JEDEC standard assumes that physically addressable media is synchronous, heavily parallel, reliable and implements a design structure that is known to a processor memory controller. Consequently, JEDEC uses a series of distinct commands that cause the DRAM devices to execute known operations in hardware.\nRecent significant development of flash-based devices enable use of non-volatile memory as a main memory replacement. However, typical non-volatile main memory solutions continue to provide communication between the processor and the non-volatile main memory using a low-level wire protocol such as JEDEC."}
{"text": "1. Field of the Invention\nThe present invention relates generally to systems and methods for absorbing energy that is generated during a collision. More particularly, the present invention involves bumper systems for vehicles wherein the structure that is used to mount the bumper to the vehicle is also designed to absorb impact energy during a crash.\n2. Description of Related Art\nBumper systems are commonly used on a wide variety of structures to protect the structure and/or its occupants during a collision. Bumper systems have been an integral and important part of automobile design for many years. As a result, many different types of systems have been developed. A popular type of bumper system is one that includes a bumper or impact member that is mounted to the vehicle using some type of mounting system that is able to absorb some of the impact energy during a collision. Mounting systems using pneumatic shock absorbers provide good absorption of impact energy and are able to withstand multiple minor impacts. However, such systems are relatively expensive and may require periodic maintenance.\nBumper supports that utilize metallic tubes which undergo lengthwise crushing during impact have also been popular. The metallic tubes provide a strong connection between the bumper and the vehicle. However, the specific energy absorption (kj/kg) of metallic tubes during crushing is not particularly high. In addition, the initial force required to initiate longitudinal crushing of a metallic tube may be too high for many situations.\nTubes made from composite materials are known to have a higher specific energy absorption than their metallic counterparts. This higher level of energy absorption is due in part to the energy absorbed when the composite tube undergoes inter-laminar splitting during longitudinal crushing. The basic problem with using composite tubes as the energy absorbing element is that the destruction of the tube that results during crushing renders the tube unable to support the bumper after a crash. In contrast, metallic tubes retain sufficient strength after crushing to provide adequate post crash support of the bumper.\nAs is apparent from the above, there is a present need to provide a simple energy absorbing system that has both the energy absorbing capabilities of a composite tube and the post-crash support strength of a metallic mounting system.\nIn accordance with the present invention, an energy absorbing composite tube is provided that is an efficient impact energy absorber and which retains good post-impact strength for maintaining a secure bumper-vehicle connection. The energy absorbing composite tube is particularly well-suited for connecting bumpers to automobiles. However, the energy absorbing composite tube may be used in any situation where two or more impact members are connected together. Typically, the invention is useful where a structure (first impact member) is protected by a bumper (second impact member) that is connected to the structure via an energy absorbing mounting system.\nEnergy absorbing composite tubes in accordance with the present invention are designed to provide an energy absorbing mount between a first impact member and a second impact member. The energy absorbing composite tubes include a mounting portion that is integral with a main tubular body portion. The mounting portion has a first impact member mounting end that is insertable into a sleeve or receptacle located in the first impact member. The mounting portion is made up of fibers embedded in a resin matrix and has an outer perimeter that is sized and shaped to matingly fit with the receptacle. The energy absorbing composite tube further includes a main tubular body portion that is integral with the mounting portion. The main tubular body portion is also made from one or more layers of fibers which are embedded in a resin matrix. The fiber layers and resin matrix may be the same or different from the mounting portion. The main tubular body portion has a perimeter that is larger than the mounting portion to provide a shoulder that extends around the perimeter of the mounting end where it meets the main tubular body.\nAs a feature of the present invention, a delamination wedge is provided that is located adjacent to the shoulder on the energy absorbing composite tube. The delamination wedge is forced against the shoulder during an impact. The resulting delamination of the main tubular body absorbs energy from the impact. The use of a delamination wedge reduces or eliminates the initial peak load that is usually necessary to initiate the delamination process. In addition, the delamination wedge eliminates the debris wedge that may form during longitudinal crushing of a tube made from composite materials.\nAs a further feature of the present invention, the mounting portion of the composite tube is not crushed and remains intact during the delamination process. The sleeve or receptacle in the impact member merely slides over the mounting portion as the delamination wedge is forced into the shoulder. The mounting portion and main tubular body retain sufficient structural integrity and strength to provide adequate post-crash support between the first and second impact members.\nThe present invention is directed not only to the energy absorbing composite tube, but also to the bumper systems and vehicles to which the energy absorbing composite tube is attached. In addition, the invention covers methods for making the energy absorbers and absorbing energy during an impact wherein the energy absorbing composite tube is utilized.\nThe present invention provides an especially useful energy absorbing mounting system that can be fine-tuned to provide a wide variety of energy absorption profiles. The number and/or types of laminates in the mounting end and main tubular body may be varied to provide a wide variety of shoulder sizes and shapes that have a direct relation to the amount of energy absorbed during delamination. In addition, the laminate layers in the main tubular body can be staggered to provide a stepped shoulder where the delamination wedge contacts each layer sequentially during a crash. Such stepped shoulder arrangements allow one to fine-tune the energy absorption profile. In addition a second mounting portion can be included at the end of the main tubular body that is mounted to the second impact member. This produces a second shoulder adjacent to the second impact member that can be used in combination with a second delamination wedge to provide even further energy absorption capabilities.\nThe energy absorbing composite tube of the present invention is a suitable replacement for metallic tube-type bumper mounting systems since it provides post-crash mounting strength equivalent to metallic tubes while at the same time providing for particularly effective energy absorption when the bumper is impacted.\nThe above described and many other features and attendant advantages of the present invention will become better understood by reference to the following detailed description when taken in conjunction with the accompanying drawings."}
{"text": "The present disclosure relates to subject matter contained in Japanese patent application No. 3-266240 (filed on Oct. 15, 1991), which is expressly incorporated herein by reference in its entirety.\n1. Field of the Invention\nThe present invention relates to a bendable portion which is provided at the distal end of an insert part of an endoscope and bent by remote control.\n2. Description of the Prior Art\nTo sterilize or disinfect an endoscope by using ethylene oxide gas or by autoclaving, it is necessary to raise and reduce the pressure in the tank containing the endoscope. Accordingly, the endoscope must be designed to withstand such a pressure change.\nThe most serious problem that is associated with autoclaving is that a bendable portion covering tube, that is the most flexible of the sheathing members of the endoscope, is inflated and likely to rupture when the pressure in the tank is reduced.\nIt has been conventional practice to provide a bellows which inflates when the ambient pressure is reduced, or to provide a valve or other similar device in the control part or the light guide connector of the endoscope so as to keep the inside of the endoscope at the same pressure as the ambient pressure, thereby preventing the bendable portion covering tube from rupturing when the pressure in the tank is reduced.\nHowever, the addition of a bellows or a valve device increases the cost of the apparatus. In addition, since the bendable portion covering structure itself is not changed, if the function of the bellows or the valve cannot satisfactorily comply with the reduced pressure conditions, the bendable portion covering tube may rupture."}
{"text": "i) Field of the Invention\nThe present invention relates to an automatic control system for lighting projectors which serve to light up a theater, studio, hall or the like for purposes of stage effects.\nii) Description of the Background Art\nIn stage lighting for a theater, it is generally necessary to topically light up a specific place on a stage by a lighting projector, for example, a spotlight or the like for the purpose of realizing desired stage effects. If an object to be lit up on the stage is a moving one such as a person or the like, or a specific place is lit up on the basis of such a moving object, it is necessary to control the lighting projector according to the movement of the object so as to shift its lighting direction.\nUsually, the lighting direction of such a lighting projector has heretofore been controlled by an operator of the lighting projector while visually observing the movement of the object.\nIn the control of the lighting direction of the lighting projector by the operator, however, the operator is required to be skilled in operating technique, and it is very difficult to reliably conduct the intended control of the lighting direction for a long period of time because of difference between individual operators in technical skill, operator's fatigue and the like. In addition, a plurality of lighting projectors are often used in lighting of this sort, and if so, the same number of operators as the lighting projectors are required.\nIn Japanese Patent Application Laid-Open No. 33802/1989, there has hitherto been proposed an automatic tracking apparatus for a spotlight, which serves to track an moving object for lighting.\nThis apparatus includes a supersonic wave transmitter held by a moving object to be lit up and adapted to reliably generate a supersonic wave at regular intervals by a clock oscillator, a plurality of supersonic wave sensors for detecting this supersonic wave, another clock oscillator in a control unit, which transmits a wave having a frequency closely similar to that of the clock oscillator held by the object, a time difference detector for detecting a time difference between the clock signal from the clock oscillator in the control unit and the clock signal from the supersonic wave sensor, a computing element for operating a distance between the supersonic wave sensor and the object on the basis of the time difference signal to locate the position of the object, and a control unit for controlling the lighting direction of the spotlight according to an output of the computing element.\nAccording to this apparatus, positional data of the moving object can be obtained theoretically. Therefore, the lighting projector can be controlled to automatically shift the lighting direction thereof, thereby tracking the moving object to light up it.\nIn practice, this apparatus however involves the following problems:\n(1) The frequencies in the supersonic wave transmitter held by the moving object and the clock oscillator in the control unit must be caused to approximate to each other with high precision. In practice, however, it is considerably difficult to achieve this approximation. This difficulty becomes greater as the number of moving objects increases. The solution of this problem requires additional signal processing means such as synchronization of a plurality of clock oscillators by radio wave signals. PA1 (2) When the apparatus is applied to a plurality of moving objects, it is necessary to use supersonic waves having frequencies separately preset on the moving objects and correspondingly, employ supersonic wave sensors corresponding to the preset frequencies. In this case, additional signal processing means are required for the individual frequencies. Therefore, the above-described difficulty becomes more and more marked. PA1 (3) The necessity of the above-described additional signal processing requires complicated signal processing for obtaining a positional data for each moving object and hence prolongs the time required for the control. As a result, the time density of the resulting positional data inevitably becomes low, leading to loose control of the lighting projector in the end. PA1 a central control unit; PA1 a radio wave transmitter for transmitting a radio wave signal according to a transmission command signal from the central control unit; PA1 a supersonic wave transmitter held by a moving object and adapted to transmit a supersonic wave by detecting the radio wave signal; PA1 a plurality of supersonic wave detecting devices each having a time counter which is initialized by a start command signal from the central control unit to start instrumentation, adapted to stop the instrumentation of the time counter by detecting the supersonic wave transmitted from the supersonic wave transmitter, and disposed in different stationary positions; and PA1 a lighting projector provided in a stationary position and having a drive mechanism for shifting its lighting direction, PA1 said central control unit including: PA1 a central control unit; PA1 a radio wave transmitter for transmitting radio wave signals according to transmission command signals from the central control unit; PA1 supersonic wave transmitters separately held by two or more moving objects and adapted to transmit respective supersonic waves by detecting the radio wave signals identified correspondingly to the supersonic wave transmitters; PA1 a plurality of supersonic wave detecting devices each having a time counter which is initialized by a start command signal from the central control unit to start instrumentation, adapted to stop the instrumentation of their corresponding time counters by detecting the supersonic wave transmitted from the supersonic wave transmitter, and disposed in different stationary positions; and PA1 at least one lighting projector provided at stationary position and having a drive mechanism for shifting its lighting direction, PA1 said central control unit including:\nAs described above, the conventional control systems for lighting projectors require complicated constitution and signal processing and consequently involve a problem that control of the lighting projectors according to the movement of moving objects cannot be achieved with high precision."}
{"text": "When developing new chip designs, it is common for chip designers to combine pre-designed components to form the new chip. This modular chip fabrication process is beneficial as it takes advantage of preexisting proven technology. For example, a chip designer might combine components from several different previously developed chips.\nThe computer chip industry is constantly developing and employing process technologies to produce chips having smaller feature sizes. For example, the achievable feature size of about three micrometers (μm) (or 3,000 nanometers (nm)) in 1976 was reduced to about 90 nm in 2003. Smaller feature sizes allow for a greater number of functionalities to be associated with a given chip, and thus generations of chip scaling have followed scaling laws first set forth by R. Dennard et al., “Design of ion-implanted MOSFETs with very small physical dimensions,” IEEE Journal of Solid State Circuits, vol. SC-9, no. 5, pp. 256-268 (October 1974). Further, chips with smaller feature sizes require less power to operate. With an estimated $2.7 billion spent in 2005 to run servers and other associated computer equipment in the United States alone, power consumption is a growing concern.\nThese rapid advances in chip scale technology, however, can have notable drawbacks with regard to modular chip fabrication. Namely, as chip technology changes, incompatibilities can often arise between the ‘old’ and the ‘new’ technologies. By way of example only, a change in chip feature size typically coincides with a change in power requirements. Different power requirements can render one component incompatible with another component. Thus, the versatility of current chip design technology can be limited unless all of the components are redesigned into the latest node for a new chip and fabricated using the scaled down features associated with the latest node of semiconductor wafers.\nTherefore, modular chip fabrication techniques which improve the compatibility of different chip technologies would be desirable."}
{"text": "1. Technical Field\nThe present invention relates to a switch device suitable for starting a vehicle engine.\n2. Related Art\nRecently a type of vehicle, in which a user does not conventionally insert a key in a key cylinder to turn the key, but the user having a proper electronic key starts up an engine only by pressing a push button of an engine starting switch device provided on a driver seat on a condition that the vehicle is equipped with an authentication system such as a so-called immobilizer, has become widespread in vehicles such as a four-wheeled vehicle. Japanese Unexamined Patent Publication No. 10-205183 discloses an automotive key cylinder in which a drain property is considered. In the automotive key cylinder disclosed in Japanese Unexamined Patent Publication No. 10-205183, a drain hole is made in a lower portion on a front-end side of a case, and a liquid (such as rain water) invading in a cylinder head from a key plate hole is drained away from the drainage hole to the outside of the case."}
{"text": "1. Field of the Invention\nThe embodiments of the invention generally relate to a slit valve door for sealing substrate passages in vacuum processing systems.\n2. Background of the Related Art\nThin film transistors (TFT) are commonly used for active matrix displays such as computer and television monitors, cell phone displays, personal digital assistants (PDAs), and an increasing number of other devices. Generally, flat panels comprise two glass plates having a layer of liquid crystal materials sandwiched therebetween. At least one of the glass plates includes one conductive film disposed thereon that is coupled to a power source. Power, supplied to the conductive film from the power source, changes the orientation of the crystal material, creating a pattern display.\nWith the marketplace's acceptance of flat panel technology, the demand for larger displays, increased production and lower manufacturing costs have driven equipment manufacturers to develop new systems that accommodate larger size glass substrates for flat panel display fabricators. Current glass substrate processing equipment is generally configured to accommodate substrates up to about five square meters. Processing equipment configured to accommodate substrate sizes exceeding five square meters is envisioned in the immediate future.\nGlass substrate processing is typically performed in a cluster tool by subjecting a substrate to a plurality of sequential processes to create devices, conductors, and insulators on the substrate. Each of these processes is generally performed in a process chamber configured to perform a single step of the production process. In order to efficiently complete the entire sequence of processing steps, the cluster tool includes a number of process chambers coupled to a central transfer chamber. A robot is housed in the transfer chamber to facilitate transfer of the substrate between the process chambers and a load lock chamber. The load lock chamber allows substrates to be transferred between the vacuum environment of the cluster tool and an ambient environment of a factory interface. Such cluster tools for glass substrate processing are available from AKT, Inc., a wholly-owned subsidiary of Applied Materials, Inc., of Santa Clara, Calif.\nAs the substrate size for manufacturing flat panel display grows, the manufacturing equipment for these substrates becomes larger in size as well. Accordingly, the door or gate that isolates one vacuum chamber (or load lock chamber) from another becomes larger, or, specifically longer, since the slot opening between the two chambers has to become wider to accommodate the large width of the substrate passing through the slot opening. The increasing length of the door poses technical challenges for obtaining a good isolation seal between the two chambers, which is maintained by an elastomer seal disposed around the slot opening between the door and a chamber wall.\nFIG. 1A depicts a partial sectional view of a substrate passage 108 formed through a chamber body 106 and selectively sealed by a conventional slit valve door 110. Conventional slit valve doors are typically comprised of a flat member of aluminum having a long lateral span. A closing force is applied toward the center of the door 110 by brackets 102 attached, as shown in FIGS. 1A-B, to a stiff rotating shaft 104. The door 110 is rotated between a position sealing the passage 108 (as shown in FIG. 1A) and a position clear of the passage 108 by an actuator 118 coupled to the shaft 104. A seal 116 is disposed between the door 110 and chamber body 106.\nThe force required to load the seal 116 in order to obtain good chamber isolation is high. The high load applied near the center of the door 110 results in a high loading force approximate the center of the door 110 and a substantially lower sealing force near the ends of the door, as depicted by force arrows 112. The shaft 104 may deflect while under load as shown by the phantom shaft 120, as the door 110 has a long span between its bearing supports 114 disposed in the walls of the chamber body 106 and the brackets 102 coupled to the center of the door 110. Deflection of the shaft 104 while the door 110 is in a closed position further aggravates the low loading condition of the seal at the ends of the door. The low sealing force at the edge of the door may lead to undesirable leakage through the passage 108.\nIn order to provide a stiffer door for more uniform seal loading, the door and/or the shaft may be fabricated from thicker materials or materials having higher modulus. However, this approach increases the cost of the load lock chamber, as high strength materials are typically expensive, and a larger load lock chamber may be required to accommodate the larger, high strength door with adequate clearance during operation. A larger load lock chamber is undesirable due to the increased material and manufacturing costs of the chamber itself, along with increased pump capacity required to pump down the larger load lock volume. Moreover, increased load lock volume typically requires increased pump time which has an adverse affect on system throughput.\nThe use of the curved slit valve has been proposed to address these concerns and is described in commonly assigned and previously incorporated U.S. patent application Ser. No. 10/867,100, entitled “CURVED SLIT VALVE DOOR”, filed Jun. 14, 2004. The implementation of a curved slit valve door has presented new engineering challenges. For example, as the door sealing surface becomes flat when pressed against the planar chamber wall to seal the slit valve passage, the change in the projected length of the curved slit valve door should be accommodated to prevent excess wear of the door actuation mechanism.\nTherefore, there is a need for an improved slit valve door."}
{"text": "Current and prior designs of undercarriage structures for track-type vehicles use design and manufacturing practices which require that the individual left and right track assemblies be parallel with each other. Tight manufacturing and assembly tolerances ensure that the left and right track assemblies are held parallel. This is done in the belief that excessive wear of the track guiding components will occur if the track assemblies are not maintained in parallel relationship. However, with the left and right track assemblies held parallel, the track links bear against the mating track roller treads in a very precise location. As the mating surfaces of the links and rollers wear, the contacting surfaces assume wear profiles which exactly match each other. Therefore, continued operation of the moving undercarriage structure produces wear along the entire contacting surfaces of the links and rollers.\nThe present invention is directed to overcoming one or more of the problems as set forth above."}
{"text": "The manufacture of silicone products generates residue that can present serious problems in its safe and environmentally acceptable disposal. A variety of methods are known for treating chlorosilane direct process residue. However, there is a persistent high level of chloride in the treated chlorosilane residue.\nMethods for reducing the level of chloride in chlorosilane residue are also known. However, the chloride level of treated residue is still too high limiting or preventing the use of residue in cement kilns. Additionally, smelter operations impose financial penalties for chlorosilane residues with chloride levels above 0.1%."}
{"text": "The present invention relates to lubrication systems for hermetic compressors and in particular vertically oriented hermetic compressors having a pick-up tube or shaft end which extends into the oil contained in a sump.\nHermetic compressors having hollow shafts or pick-up tubes which extend into the oil contained in a sump are used for drawing lubricant from the sump and providing it to moving parts of the compressor assembly. Often these tubes or shafts are provided with a paddle for pumping the oil through the shaft or tube. A problem experienced with pick-up tubes or shafts is that on their rotation they tend to create vortices around the end of the shaft or tube which extends into the oil. These vortices hinder the performance of oil delivery into the tube or shaft. A way of preventing the formation of these vortices is desirable."}
{"text": "1. Field of the Invention\nThe invention relates to systems used for chemical sterilization of medical devices, and more particularly, to systems having multiple chambers used for chemical sterilization of medical devices.\n2. Description of the Related Art\nMedical instruments have traditionally been sterilized using either heat, such as is provided by steam, or a chemical, in the gas or vapor state. Sterilization using hydrogen peroxide vapor has been shown to have some advantages over other chemical sterilization processes.\nThe combination of hydrogen peroxide with a plasma provides certain additional advantages, as disclosed in U.S. Pat. No. 4,643,876, issued Feb. 17, 1987 to Jacobs et al. U.S. Pat. No. 4,756,882, issued Jul. 12, 1988 also to Jacobs et al. discloses the use of hydrogen peroxide vapor, generated from an aqueous solution of hydrogen peroxide, as a precursor of the reactive species generated by a plasma generator. The combination of hydrogen peroxide vapor diffusing into close proximity with the article to be sterilized and plasma acts to sterilize the articles and remove residual hydrogen peroxide. However, effective sterilization of articles having long narrow lumens are very difficult to achieve, since the methods are dependent upon diffusion of the sterilant vapor into close proximity with the article before sterilization can be achieved. Thus, these methods have been found to require high concentration of sterilant, extended exposure time and/or elevated temperatures when used on long, narrow lumens. For example, lumens longer than 27 cm and/or having an internal diameter of less than 0.3 cm have been particularly difficult to sterilize. The sterilization of articles containing long narrow lumens therefore presents a special challenge.\nU.S. Pat. No. 4,744,951 to Cummings et al. discloses a two-chambered system which provides hydrogen peroxide in vapor form for use in sterilization processes. The sterilant is initially vaporized in one chamber and then applied to the object to be sanitized in another single sterilizing chamber, thereby producing a concentrated hydrogen peroxide vapor which is relatively more effective. The sterilization processes are designed for furnishing concentrated hydrogen peroxide vapor to interior surfaces of articles having a tortuous or a narrow path. However, the sterilization processes are ineffective at rapidly sterilizing lumened devices, since they depend on the diffusion of the hydrogen peroxide vapor into the lumen to effect sterilization.\nU.S. Pat. No. 4,797,255 to Hatanaka et al. discloses a two-chambered sterilization and filling system consisting of a single sterilization chamber adjacent to a germ-free chamber utilized for drying and filling sterilized containers.\nU.S. Pat. No. 4,863,688 to Schmidt et al. discloses a sterilization system consisting of a liquid hydrogen peroxide vaporization chamber and an enclosure for sterilization. The enclosure additionally may hold containers wherein the hydrogen peroxide sterilant vapor does not contact the interior of the containers. This system is designed for controlling the exposure to the hydrogen peroxide vapor. The system is not designed for sterilizing a lumen device.\nU.S. Pat. No. 4,952,370 to Cummings et al. discloses a sterilization process wherein aqueous hydrogen peroxide vapor is first condensed on the article to be sterilized, and then a source of vacuum is applied to the sterilization chamber to evaporate the water and hydrogen peroxide from the article. This method is suitable to sterilize surfaces, however, it is ineffective at rapidly sterilizing lumened devices, since it too depends on the diffusion of the hydrogen peroxide vapor into the lumen to effect sterilization.\nU.S. Pat. No. 4,943,414, entitled “Method for Vapor Sterilization of Articles Having Lumens,” and issued to Jacobs et al., discloses a process in which a vessel containing a small amount of a vaporizable liquid sterilant solution is attached to a lumen, and the sterilant vaporizes and flows directly into the lumen of the article as the pressure is reduced during the sterilization cycle. This system has the advantage that the water and hydrogen peroxide vapor are pulled through the lumen by the pressure differential that exists, increasing the sterilization rate for lumens, but it has the disadvantage that the vessel needs to be attached to each lumen to be sterilized.\nU.S. Pat. Nos. 4,937,046, 5,118,471 and 5,227,132 to Anderson et al. each disclose a sterilization system which uses ethylene oxide gas for sanitation purposes. The gas is initially in a small first enclosure and thereafter slowly permeates into a second enclosure where the objects to be sterilized are located. A medium is then introduced into the second enclosure to flush out the sterilizing gas into a third enclosure containing the second enclosure. An exhaust system then exhausts the sterilant gas and air from the third enclosure. These systems also have the disadvantage of relying on the diffusion of the sterilant vapor to effect sterilization and hence are not suitable for rapidly sterilizing lumened devices.\nU.S. Pat. No. 5,122,344 to Schmoegner discloses a chemical sterilizer system for sterilizing items by vaporizing a liquid chemical sterilant in a sterilizing chamber. Pre-evacuation of the sterilizer chamber enhances the sterilizing activity. Sterilant is injected into the sterilizer chamber from a second prefilled shot chamber. This system also relies upon diffusion of sterilant vapor to effect sterilization and is also not suitable for rapidly sterilizing lumened devices.\nU.S. Pat. No. 5,266,275 to Faddis discloses a sterilization system for disinfecting instruments. The sterilization system contains a primary sterilization chamber and a secondary safety chamber. The secondary safety chamber provides for sensing and venting to a destruction chamber any sterilization agent that is released from the primary sterilization chamber. This system, as in other systems, also relies upon diffusion of sterilant vapor to effect sterilization and is also not suitable for rapidly sterilizing lumened devices.\nIn U.S. Pat. Nos. 5,492,672 and 5,556,607 to Childers et al, there is disclosed a process and apparatus respectively for sterilizing narrow lumens. This process and apparatus uses a multicomponent sterilant vapor and requires successive alternating periods of flow of sterilant vapor and discontinuance of such flow. A complex apparatus is used to accomplish the method. Additionally, the process and apparatus of '672 and '607 require maintaining the pressure in the sterilization chamber at a predetermined subatmospheric pressure.\nIn U.S. Pat. No. 5,527,508 to Childers et al., a method of enhancing the penetration of low vapor pressure chemical vapor sterilants into the apertures and openings of complex objects is disclosed. The method repeatedly introduces air or an inert gas into the closed sterilization chamber in an amount effective to raise the pressure to a subatmospheric pressure to drive the diffused sterilant vapor further into the article to achieve sterilization. The '508, '672 and '607 Childers inventions are similar in that all three require repeated pulsations of sterilant vapor flow and maintenance of the sterilization chamber pressure at a predetermined subatmospheric pressure.\nIn U.S. Pat. No. 5,534,221 to Hillebrenner et al., a device and method for sterilizing and storing an endoscope or other lumened medical device is disclosed. The device includes a sealable cassette in which the endoscope or other medical device is placed. The cassette has an input port for receiving a sterilizing agent through a connector, an output port for expelling the sterilizing agent when a vacuum is applied thereto through a connector, and check valves in the input and output ports to open the ports when the connectors are coupled to the ports and to seal the ports when the connectors are removed from the ports such that after the endoscope has been sterilized, it remains sterilized within the cassette until the cassette is opened. The method of the '221 invention involves placing the medical device inside the cassette and coupling the device to either the input or output port of the cassette. The cassette is then placed in an outer oven-like container or warming chamber where the temperature is properly maintained. Connections are made to open the input and output ports on the cassette such that the sterilizing agent may be introduced through a first port to bathe the outside of the medical instrument or other object, such as an endoscope while one end of the hollow object, such as the endoscope, is coupled to the output port where a vacuum is supplied external to the cassette to pull the sterilization agent into the cassette and through the interior passageways of the endoscope. When the sterilization process is completed, the warming chamber is opened and the sterilizing cassette is simply removed from the chamber with the input and output ports being uncoupled from their respective sources. A tight seal is maintained and the object remains in the sterilized interior of the cassette until the cassette is opened or the device is to be used. Thus, the '221 invention is concerned with providing a means whereby a sterilized medical device can be retained within a cassette in which it was sterilized until ready for use, thus avoiding any contamination by exposure to the atmosphere or handling before use. Additionally, in some cases of the '221 invention, wherein the lumen of the device to be sterilized is connected to the output port, particularly wherein the devices have long, narrow lumens, the time to expel the sterilizing agent through the lumen and out of the cassette may be undesirably long. Also, in cases wherein the lumen device is very flexible, lumen collapse may occur, either slowing or preventing vapor exit or causing lumen damage.\nU.S. Pat. Nos. 5,445,792 and 5,508,009 to Rickloff et al. each disclose a sterilization system essentially equivalent to the system disclosed in Hillebrenner '221.\nU.S. Pat. No. 5,443,801 to Langford teaches a transportable cleaning/sterilizing apparatus and a method for inside-outside sterilization of medical/dental instruments. The apparatus avoids the use of heat, pressure, severe agitation, or corrosive chemicals which might damage delicate equipment. This invention uses ozone gas or solution as sterilant. It does not involve the use of sterilant vapor or vaporizing a sterilant solution into vapor, and is not suitable for operations under vacuum because flexible bags or containers are used.\nIn consideration of the foregoing, no simple, safe, effective method of sterilizing smaller lumens exists in the prior art. Thus, there remains a need for a simple and effective method of vapor sterilization of articles with both long, narrow lumens as well as shorter, wider lumens. Furthermore, there also remains a need for a simple and effective sterilization system with independently operable chambers."}
{"text": "Computer based document handling systems are generally divided into four broad categories: text editors and word processing systems; formatters; syntax directed editors; and specialized tools. Most systems have features from more than one of these four broad categories in addition to any image processing necessary for editing images. An image processor reads an image of a document using an image input unit, stores it in a memory in the form of image data, and subjects the stored image data to editing operations such as addition and deletion. As one of the editing operations, part of the read document is electronically cut and pasted to a predetermined place in another document. In such an electronic cutting and pasting operation, the document image to be cut and the destination document image are displayed either on a single display unit or on separate display units, and a cutting image portion in the document to be cut and a pasting place in the pasting (destination) document are similarly specified by a mouse, for example. In the related art image processor, the cutting image portion and the pasting place are specified by the operator by selecting appropriate positions while looking at the screen. For this reason, it is not possible to cut or paste figures precisely. Particularly, shifts in position are noticeable when the same patterns must be juxtaposed. U.S. Pat. No. 5,224,181 to Tsutsumi entitled “Image processor” and U.S. Pat. No. 5,202,670 to Oha entitled “Image Processing Apparatus” both discuss image processing generally.\nFor documents in which presentational considerations are important, the documents must be submitted to a formatter for preparation prior to presentation. Formatters are non-interactive tools that process a document to produce either a display independent or a device dependent layout specification. Documents are submitted to formatters in the form of descriptions on file and they carry out the processing and return the overall results after a certain period of time. High-level formatters work on the basis of a logical description of the document. The user is not required to specify the presentation details desired. The user deals with the logical organization of the document, i.e., the different types of elements that appear in the document, such as, for example, section, paragraph, heading, summary, etc. The formatter handles the layout presentation of these elements. Low-level formatters make it possible to include commands within the document description to enable changes in other characteristics of the document, such as, for example, font, spacing, margins and justification. U.S. Pat. No. 5,438,512 to Mantha et al. entitled “Method and apparatus for specifying layout processing of structured documents” is directed to high-level formatters. Mantha discloses a method and apparatus for specifying layout processing of logically structured documents in computer document handling systems. The Mantha method and apparatus allow the specification of the generic logical structure of the structured document in terms of relational attribute grammars.\nMost interactive systems allow the user to see the layout of a document as it is being prepared. These interactive systems also separate the logical structure of the documents from the specification of the presentation details. Typically, interactive systems as well as high-level formatters, use a grammatical notation to describe the logical structure of documents. These logical structures are mostly hierarchical in nature and tree structures are used to represent them.\nThe need for document processing devices which can generate a layout structure of a document by applying a template (i.e., rules for layout) to a logical structure of the document are apparent with the wide use of microprocessors to process and properly render electronically created or digitally scanned documents, or images as more specifically referred to at times. In general, both of the logical development and the layout of an actual output document are important factors of the document. However, the layout is not necessarily important in the process of producing a document. For example, although the layout is necessary once the contents of a document are determined, in some cases the layout is not considered in the initial stage of document production where the logical development is not clear yet. Furthermore, different layouts may be required for one logical development in some cases; for instance, when the same document needs to be distributed to a plurality of persons or sections. Recognizing the above distinction, there have been proposed techniques of generating a layout representing such structures called “layout structure” from a structure called “logical structure” which represents the chapter construction of a document. The process of generating a layout structure from a logical structure is called a layout process, and is performed by a program (layout processing program) that is incorporated in document processing systems.\nIn order to generate various forms of layout structures from the same logical structure, the operation of a layout processing program needs to be modified. However, in general, a program itself cannot be modified properly by ordinary users because the modification needs expert knowledge. For this reason, there has been employed a technique of altering a program by using parameters such as “page size is A4” and “double column setting” for textual documents. However, according to this technique, the kinds of specifiable parameters are limited to ones that are preset in the program. Furthermore, to control the operation of a layout processing program having many functions it is necessary to specify a large number of parameters which precludes ordinary users from using such a program. To solve the above problems, there have been proposed techniques of controlling a layout processing program by using, rather than simple parameters, a data structure representing a layout template. One of those techniques is a “generic layout structure” prescribed in the international standards “ODA” (ISO8613, Information Processing-Text and Office Systems-Office Document Architecture (ODA) and Interchange Format (1989)). The ODA only sets forth data structures for representing document structures and guidelines for their use, and does not describe actual layout processes. However, it is apparent that the following functions are needed to perform layout operations according to the ODA. The layout process having the following functions is hereinafter referred to as “ODA layout process.”\n(1) Layout processing function based on a layout template (generic layout structure)\n(2) Layout process selecting function (top-down or bottom-up)\n(3) Reuse of a layout result\n(4) Category-based layout processing function\nThe functions (mechanisms) (1)-(4) are not necessarily effected individually, but could be combined when desired. Combining the functions can improve the efficiency of the layout process.\nU.S. Pat. No. 5,381,523 issued to Hayashi entitled “Document processing device using partial layout templates” discloses partial layout templates that are prepared for respective partial logical structures of a hierarchical logical structure of a document. Each partial layout template expresses rules for producing a layout of the corresponding partial logical structure. A partial layout generator produces a partial layout structure by recursively calling itself or by calling a content layout system while referring to the partial layout template, and pours the generated partial layout structure into a lowest-rank frame. A layout of the entire logical structure is produced by sequentially performing the partial layout operations. Hayashi also only relates to a document processing device which can generate a layout structure of a single document by applying a template (i.e., rules for layout) to a logical structure of the single document.\nWith the growing interest in digital photography and the necessity to lay out and view several digital images as a document, a need still exists for a document/image structuring process which would allow a microprocessor to organize the layout of several images on a sheet-like medium. Needed is a technology that enables a user to gang scan (several images on the scanner platen at once) many images and automatically locate each image, crop them, and correct any rotation errors associated with each image. Resulting images can then be stored individually or as a structured image with a user-defined or simple row-column layout which can be created. A detailed description of Structured Images is provided in U.S. Pat. No. 5,485,568 issued to Venable et al. on Jan. 16, 1996, entitled “Structured Image (SI) format for describing complex color raster images.”\nThere is a need for technology that will automatically generate the digital equivalent of a photographers contact print sheet in which the segmented images will be scaled and distributed on a resulting page such that each image is scaled as large as possible wherein there is a minimum of white space remaining on the rendered page.\nAs a first constraint to adequately addressing such a need, all images must be scaled by the same amount, i.e., the relative size of each segmented image stays the same. A second constraint should be in forming the equivalent of a bounding box with minimum white space (space uncovered by an rectangle). Once a solution to the above constraints are found, many images may be scaled by the appropriate amount to make a bounding box the same size as the printable area of paper. It is a feature of the present invention to effectively address the constraints in order to accomplish the desired results described above.\nOther advantages and salient features of the invention will become apparent from the detailed description which, taken in conjunction with the drawings, disclose the preferred embodiments of the invention."}
{"text": "In a batch type hot rolling line in which a steel piece is individually heated and rolled by a roughing mill and a finishing mill to make a steel sheet of a desired thickness, a line stop tends to occur because the rolled material does not successfully bite into a gap defined by the upper and lower rolls and faulty shapes of the leading and the trailing ends of the rolled material will considerably lower the product yield.\nTherefore, an improved rolling method (endless rolling method) has been adopted in recent years in which the trailing end of a steel piece to be rolled is joined to the leading end of the succeeding steel piece to be rolled before the finish rolling and the resulting steel piece is continuously supplied to a hot rolling line.\nIn this respect, reference may be had to Japanese Unexamined Patent Publication No. 58-122,109 which describes a method wherein the trailing end of a preceding steel piece is butt-joined to the leading end of a succeeding steel piece over the entire area of the end plane thereof, before they are rolled. Also, Japanese Unexamined Patent Publication No. 4-89,120 describes a method in which the trailing end of a preceding steel piece is arranged against the leading end of a succeeding steel piece with a gap on the entry side of a hot finishing mill equipment and in this area an alternate magnetic field is applied to the steel pieces in their thickness direction to heat them and raise the temperatures thereof by induction heating and then both steel pieces are joined by pressing.\nIn such continuous hot rolling of steel pieces, a cutting step is added as a preparatory step for the joining process, in which defective portions at the leading and the trailing ends of the steel piece are cut by a crop shear, drum shear or the like. In the cutting step, however, there may be instances wherein the leading end and/or the trailing end of the steel piece tend to warp upward or downward. As a result, the steel pieces may be joined with an upward or downward displacement (hereinafter referred to as \"displacement\"), as they are pressed against each other.\nIn this respect, the reason why the end of the steel piece warps in the cutting step is that a moment is produced in the cut plane when the end of the steel piece is cut by a crop shear. Usually, the preceding steel piece and the succeeding steel piece are cut by the same cutting blade so that they warp in opposite directions. If the steel pieces are heated to raise their temperatures and pressed in such a condition, they tend to further warp upward and downward to substantially decrease the effective joint area. This tends to cause the demerit that the steel pieces are broken in the following finish rolling process at the joined portion. In particular, when a transverse-type induction heating is used as heating means and an alternate magnetic field is applied to the steel pieces in their thickness direction for heating them by induction heating, due to the requirement for the installation space of the heating coil, a clamp for sandwiching the steel pieces must be used at a location which is remote from the ends of the steel pieces. As a result, it has been considered that the influence of the warp in the pressing process of the steel pieces would be more significant by using a clamp having a pair of upper and lower jaw elements which extend toward the ends of the steel pieces.\nFurthermore, Japanese Unexamined Patent Publication No. 5-185,111 describes a cutting device having a structure wherein two pairs of upper and lower drums are arranged along the transfer direction of the rolled material and respectively rotated in opposite directions from each other, and cutting blades are mounted on each pair of drums and oriented in opposite directions. Also, Japanese Unexamined Patent Publication No. 56-27,719 and Japanese Unexamined Patent Publication No. 56-119,311 each describes a cutting technique in which the shear members of a drum type shear are driven at different circumferential speeds to cut the crop end of the steel piece. Moreover, Japanese Unexamined Patent Publication No. 7-251,203 describes a cutting technique in which a flying shear for cutting the end of the steel piece in a parallel surface relation is applied when a plurality of elongate steel pieces for hot rolling are joined by laser welding and are continuously hot rolled. However, even when such device and/or techniques are applied to the cutting of the steel pieces, it is still impossible to reduce the warping of the steel piece which causes the displacement at the joined portion.\nIt is an object of the present invention to provide a method for reducing the displacement at the joined portion of the steel pieces even when a traverse-type high-frequency induction heating is used as heating method, thereby allowing a stable continuous hot rolling to be carried out without defective joining or breakage at the joined portion."}
{"text": "There are methods for manufacturing glass plates by employing down-draw process, as disclosed in Patent Literature 1 (Japanese Patent Application Laid-Open Publication JP-A-2004-115357). In down-draw process, molten glass is first poured into a forming member, and then the molten glass is made to overflow from the top sections of the forming member. The molten glass that has overflowed then flows downward along the opposite side surfaces of the forming member, and the streams of molten glass merge at the lower end of the forming member, thus being made into a sheet-form glass (sheet glass). The sheet glass is then drawn downward by rollers and cut into predetermined lengths."}
{"text": "A mass spectrometer is an instrument used to measure the mass, or more specifically the mass to charge ratio, of ionized atoms or electrically charged particles. Mass spectrometers help determine the composition of an unknown sample by isolating ionized atoms based on their mass-to-charge ratio, measured in Atomic Mass Units per charge (AMU/q). Mass spectrometers find widespread application in the basic sciences, medicine, and space-based research. Two common space-related applications of mass spectrometry are the study of the composition of planetary atmospheres and the monitoring of air quality on manned space missions. Although mass spectrometry has been used in space-related applications for many years, usage in space presents unique design challenges, both in terms of detection sensitivity and logistical considerations such as weight and power requirements.\nIn early mass spectrometers, atoms or molecules were ionized by a hot filament and accelerated through the instrument under the influence of voltage gradients. The ions followed a semi-circular trajectory through the instrument, which utilized a strong magnetic field to selectively direct ions of a specific mass towards a detector. By controlling the strength of the magnetic field and the accelerating voltage, ions of different mass/charge ratios could be selectively guided towards the detector. These early mass spectrometers suffered from numerous deficiencies and drawbacks, most significantly the difficulty in achieving and maintaining a stable magnetic field.\nQuadrupole mass spectrometers (QMS) eliminated the need for magnetic fields. Similar to its predecessor, a QMS employs a hot filament to ionize the atoms or molecules. Ionization results from the conversion of normally neutral atoms or molecules to electrically charged particles. The ions are accelerated through a mass filter having four parallel metal rods, referred to as the quadrupole. DC and RF (frequency Ω) voltages are applied to opposing pairs of these rods with opposite polarities to create an electric field inside the rod assembly. For a given DC and RF voltage, only ions of a certain mass-to-charge ratio will pass through the quadrupole filter, while all other ions are thrown out of their original path. The stability region is defined by ion trajectories that are periodic and bounded. A detector placed at the end of the rod assembly opposite the ionizer measures those ions that pass through the quadrupole filter. A mass spectrum is obtained at the detector by measuring the ions passing through the quadrupole filter as the voltages on the quadrupole rods are varied. The mass resolution of the QMS is the maximum atomic mass that can be distinguished. The maximum attainable resolution is determined by both the fidelity of the electronics and the overall tolerances of the instrument design. Generally, the voltages employed in QMS systems are of the order of a few thousand volts to obtain a mass resolution of a few hundred Daltons.\nThe quadrupole rods can be a circular or hyperbolic. Circular rods are easier to manufacture and consequently cheaper. However, the quadrupole electric field produced with circular rods is slightly distorted, which can reduce the maximum attainable mass resolution of the instrument. Consequently, in applications requiring high mass resolution, the more difficult and expensive to manufacture hyperbolic rods are employed as quadrupole rods.\nTo improve resolution, the electric field generated by the quadrupole rods can be perturbed by introducing an excitation RF signal at an auxiliary frequency (ω) different from the fundamental frequency (Ω). This perturbation causes the original stability region to break into smaller regions termed islands, including an ‘upper stability island.’ The result of this auxiliary frequency is the creation of bands of instability in the previously stable regions of the electric field. Charged particles having a mass within a certain range that previously passed through a stable region of the electric field may now be thrown off trajectory as they coincide with these islands of instability. In this way, the use of an auxiliary frequency to drive the quadrupole rods allows a QMS to operate with improved resolution. A QMS driven under an auxiliary frequency excitation is able to better differentiate between charged particles having close, yet different masses, or mass-to-charge ratios. The size and shape of the upper stability island is determined by the auxiliary frequency used and the amplitude of the excitation RF signal. To create an island of appropriate size, for example, the auxiliary frequency (ω) inserted into the QMS system needs to be near an integer multiple of the fundamental RF frequency (i.e., ω=0-0.1Ω, 0.9-1.1Ω, 1.9-2.1Ω, etc.). Employing an excitation RF signal in one of these auxiliary frequency ranges in conjunction with the DC voltage U and the RF voltage V allows for improved resolution and discrimination between ions with small differences in their mass-to-charge ratio. In general, the auxiliary signal amplitude required for appropriate island formation increases with auxiliary signal frequency.\nUnfortunately, this use of auxiliary frequency excitation presents problems in constrained applications, such as space-based applications, where it is advantageous for the QMS to have increased sensitivity and enhanced resolution to better detect and differentiate between complex molecules with higher masses. Having to excite the quadrupoles with an excitation RF signal at an auxiliary frequency in order to create islands of stability/instability requires higher power and increased complexity of the voltage control system. Additionally, the excitation RF signal must be driven at an amplitude that corresponds to a few hundred volts (˜10% of the fundamental RF signal amplitude), to create islands of the appropriate size. However, in space-based applications, power and size is at a premium.\nOther factors that affect the resolution and accuracy of the measurement made by the QMS are imperfections in the rods and limitations of the electronics. Furthermore, electronic component values drift with temperature and time, which can have the material effect of shifting the operating point of the quadrupole sufficiently to degrade the detected mass spectrum.\nThus there exists a need to enable a QMS to resolve species of heavy, complex molecules in a power efficient manner, but also to improve the tolerance of the instrument to variations in electronic component values."}
{"text": "Insect pests are a major factor in the loss of the world's agricultural crops. For example, armyworm feeding, black cutworm damage, or European corn borer damage can be economically devastating to agricultural producers. Insect pest-related crop loss from European corn borer attacks on field and sweet corn alone has reached about one billion dollars a year in damage and control expenses.\nTraditionally, the primary method for impacting insect pest populations is the application of broad-spectrum chemical insecticides. However, consumers and government regulators alike are becoming increasingly concerned with the environmental hazards associated with the production and use of synthetic chemical pesticides. Because of such concerns, regulators have banned or limited the use of some of the more hazardous pesticides. Thus, there is substantial interest in developing alternative pesticides.\nBiological control of insect pests of agricultural significance using a microbial agent, such as fungi, bacteria, or another species of insect affords an environmentally friendly and commercially attractive alternative to synthetic chemical pesticides. Generally speaking, the use of biopesticides presents a lower risk of pollution and environmental hazards, and biopesticides provide greater target specificity than is characteristic of traditional broad-spectrum chemical insecticides. In addition, biopesticides often cost less to produce and thus improve economic yield for a wide variety of crops.\nCertain species of microorganisms of the genus Bacillus are known to possess pesticidal activity against a broad range of insect pests including Lepidoptera, Diptera, Coleoptera, Hemiptera, and others. Bacillus thuringiensis (Bt) and Bacillus papilliae are among the most successful biocontrol agents discovered to date. Insect pathogenicity has also been attributed to strains of B. larvae, B. lentimorbus, B. sphaericus (Harwook, ed., ((1989) Bacillus (Plenum Press), 306) and B. cereus (WO 96/10083). Pesticidal activity appears to be concentrated in parasporal crystalline protein inclusions, although pesticidal proteins have also been isolated from the vegetative growth stage of Bacillus. Several genes encoding these pesticidal proteins have been isolated and characterized (see, for example, U.S. Pat. Nos. 5,366,892 and 5,840,868).\nMicrobial insecticides, particularly those obtained from Bacillus strains, have played an important role in agriculture as alternatives to chemical pest control. Recently, agricultural scientists have developed crop plants with enhanced insect resistance by genetically engineering crop plants to produce pesticidal proteins from Bacillus. For example, corn and cotton plants have been genetically engineered to produce pesticidal proteins isolated from strains of Bt (see, e.g., Aronson (2002) Cell Mol. Life Sci. 59(3):417-425; Schnepf et al. (1998) Microbiol Mol Biol Rev. 62(3):775-806). These genetically engineered crops are now widely used in American agriculture and have provided the farmer with an environmentally friendly alternative to traditional insect-control methods. In addition, potatoes genetically engineered to contain pesticidal Cry toxins have been sold to the American farmer. While they have proven to be very successful commercially, these genetically engineered, insect-resistant crop plants provide resistance to only a narrow range of the economically important insect pests.\nAccordingly, there remains a need for new Bt toxins with a broader range of insecticidal activity against insect pests, e.g., toxins which are active against a greater variety of insects from the order Lepidoptera. In addition, there remains a need for biopesticides having activity against a variety of insect pests and for biopesticides which have improved insecticidal activity."}
{"text": "Document management systems allow users to create centralized repositories, or libraries, containing all of the data they generate, such as information stored in documents, spreadsheets, text files, electronic mail, multimedia, etc. Powerful search and retrieval tools make this information easily available for use and collaboration across the entire enterprise. In certain instances, a user may require that a certain document or other electronic file not be widely disseminated or that the document have restricted access.\nThe widespread dissemination of electronic documents across the world wide web and other wide area networks or metropolitan area networks has complicated not only the management and integration of access control systems but also the types of access that are to be granted to particular users. For example, users from different organizations that are collaborating on a particular project may desire access to each other's electronic documents, but the access control lists and associated processes and access privileges may be incompatible and prove difficult to harmonize."}
{"text": "1. Field of the Invention\nOne disclosed aspect of the embodiments relates to an image pickup apparatus and an image pickup system.\n2. Description of the Related Art\nIn recent years, owing to further improvement in the performance of image pickup apparatuses, there has been studied a configuration including a charge holding unit within a pixel separately from a photoelectric conversion unit and a floating diffusion (hereinafter, FD). As for usage of the holding unit, the holding unit is provided to realize a global electronic shutter as disclosed in Japanese Patent Laid-Open No. 2011-216969."}
{"text": "Laboratory processes involving nucleic acid sequences are nowadays very common and often performed as a matter of routine. Such processes include inter alia hybridizing and enzymatic reactions.\nA common type of such processes is amplification reactions, such as the polymerase chain reaction, or abbreviated PCR. As is well-known, the PCR technique provides for the highly specific amplification of unique DNA segments defined by two surrounding primer sequences. PCR thereby offers a convenient way of obtaining sufficient quantities of DNA for inter alia nucleotide sequencing purposes. One major application of PCR is for diagnostic purposes. For a description of PCR it may, for example, be referred to White, T. et al., Trends in Genetics, 5, 179 (1989).\nExtensively used solid supports in the context of molecular-genetic reactions have so far been paramagnetic beads due to the large total area provided thereby and their simplified handling and processing. Thus, for example, such magnetic beads having streptavidin immobilized thereto are commercially available.\nHowever, the simultaneous processing of large sets of samples through sequential reaction steps using e.g. magnetic beads as a solid phase is technically demanding and involves a substantial risk of contamination between reactions. This is, of course, particularly undesired in amplification contexts, such as PCR, where the multiplication of contaminating sequences may result.\nDD-A-279 506 discloses the use of a multipronged device for DNA-sequencing on a solid phase by chemical degradation (Maxam-Gilbert method). The device has a set of rods, each with immobilization primers attached, which rods are designed to fit into a set of reaction vessels. The labelled DNA fragments to be sequenced are immobilized to the rods by dipping them into respective vessels containing the DNA fragments. The further processing of the immobilized DNA fragments is conducted by dipping the rods into vessels containing corresponding reagents and solution, and the contents of the respective vessels containing DNA fragments degraded by the respective base specific reagents are lyophilized and then subjected to gel electrophoresis.\nIt is readily seen that the use of a multipronged device as suggested by Rosentahl et al., supra, obviates several of the problems related to the use of separate solid phase elements, like paramagnetic beads or microtiter wells, for example. A problem of this \"patrix-matrix\" type approach, however, which is likely to have limited its application, has been to provide for sufficient binding capacities on the individual prongs.\nIt may be mentioned in this context that a similar \"patrix-matrix\" strategy, permitting sets of solid supports (patrices) to be coordinately moved between corresponding sets of reaction wells (matrices), has previously been applied in peptide synthesis. Also, a type of multipronged solid support permitting the simultaneous processing of multiple samples is commercially available for immunoassay applications.\nThus, while a nucleic acid sample procedure using a multipronged device as described above may simplify the procedure significantly and reduce the risk of mix-up and contamination to a substantial degree, there still remains, however, the relatively cumbersome step of releasing the reaction products from the solid support and transferring them to the analyzer apparatus in question."}
{"text": "This application claims the priority of Japanese Patent Application No. 2000-41679 filed Feb. 18, 2000, and Japanese Patent Application No. 2000-85972 filed Mar. 27, 2000.\n1. Field of the Invention\nThe present invention relates to an endoscope apparatus, and more particularly to an apparatus capable of changing field curvature characteristics when an object to be observed is optically enlarged among others.\n2. Description of the Prior Art\nIn recent years, in an electronic endoscope apparatus or the like, a movable lens for variable-power has been disposed at an objective lens system at the tip end portion of a scope to optically enlarge an image of the object to be observed. This optically enlarged image is picked up by an imaging device such as CCD (Charge Coupled Device), and a video signal (image signal) outputted from this CCD is subjected to various image processing by a processor apparatus to thereby display an enlarged image for the object to be observed on a monitor.\nFIGS. 16(A) and 16(B) show an example of an enlarged display image formed by the electronic endoscope apparatus, and FIG. 16(A) shows a display on the monitor, and FIG. 16(B) shows an actual cross section showing the displayed object to be observed. This FIG. 16(A) enlarges and displays a cave-in K1 of a diseased part within a circular mask M of a monitor 1, and a variable-power mechanism in the endoscope is adapted to allow a watched part to be observed by image enlargement by 70 to 100 times.\nObject of the Invention\nIn an objective lens system of the endoscope, however, when a particularly uneven part is enlarged, astigmatism, curvature of field or the like may throw a peripheral part or a central part out-of-focus. For example, a case of FIGS. 16(A) and 16(B) will be explained. When the central part of an enlarged cave-in K1 is in focus, the periphery part goes out of focus, and when the peripheral part of the cave-in K1 is in focus, the central part goes out of focus, and the entire image may not uniformly in focus. In such a state, there is the inconvenience that particularly in a freeze-frame picture to be photographed for recording, any completely satisfactory image cannot be obtained.\nAlso, in the above-described enlarged image for the uneven part, it is comparatively difficult to accurately grasp what degree the depth of the recessed part, the height of the protruded part, or the like is, and if height (difference) information on this unevenness could be provided, it would be exceedingly useful to diagnose, and to deal with the watched part among others.\nThe present invention has been achieved in the light of the above described problems, and is aimed to provide an endoscope apparatus capable of obtaining an image by attaining excellent entire focus especially during enlargement correspondingly to a shape of unevenness or the like of the object to be observed through the use of field curvature characteristics of the lens, and of obtaining information on differences in height of the uneven shape.\nIn order to attain the above described object, there is, according to the present invention, provided an endoscope apparatus including: an objective lens group disposed at a tip end portion of the endoscope, capable of changing the field curvature characteristics (image surface curvature characteristics including astigmatism) by moving a movable lens; image surface curving means for operating so as to change the field curvature characteristics of this objective lens group; and driving means for driving the movable lens on the basis of the operation of the image surface curving means.\nAccording to the present invention, there are provided two movable lenses (or lens groups) in, for example, an objective lens group, whereby the structure is arranged so as to enable both a variable-power operation and an field curvature (characteristic change) operation. Thus, when an image surface curving switch disposed in an operating unit or the like is operated, the movable lens can be driven by an actuator or the like to thereby change the field curvature characteristics. This field curvature can be caused to occur in an under direction or in an over direction with respect to a reference plane, and it becomes possible to attain focus in accordance with a shape of the recessed part by an operation in the Under direction, or in accordance with a shape of the protruded part by an operation in the Over direction.\nAccording to another aspect of the present invention, there is provided an endoscope apparatus, including: an objective lens group, disposed at the tip end portion of the endoscope, capable of changing the field curvature characteristics by moving a movable lens; driving means for driving a movable lens for field curvature of this objective lens group; judging means for dividing an imaging area for the objective lens group into a plurality and judging whether or not the central part and the peripheral part are in focus within a predetermined range by comparing information from these divided areas; and a control circuit for executing various control on the basis of output from this judging means.\nAccording to the another aspect of the present invention, it is compared, by comparing means, whether or not the central part and the peripheral part are in focus and this focusing or defocusing, that is, the focused state is supplied to the control circuit. This control circuit outputs, for example, a static image when the whole is in a focused state. That is, when a freeze switch is depressed, no static image is outputted if the central part and the peripheral part are not in focus, but the freeze-frame picture can be displayed on the monitor and recorded only when in focus. Also, this information on the focused state may be displayed on the monitor or the like, and be notified by another means.\nAccording to a further aspect of the present invention, there is provided an endoscope apparatus, including: an objective lens group, disposed at the tip end portion of the endoscope, capable of changing the field curvature characteristics by moving a movable lens; driving means for driving a movable lens for field curvature of this objective lens group; operating means for operating height (difference) information of the image central part to the peripheral part from an amount of field curvature changed by the objective lens group; and output means for outputting information on differences in height obtained by this operating means.\nIn the further aspect of the present invention, there can be provided display control means for displaying the information on differences in height obtained by the operating means on the monitor, and contour information can be displayed as the information on differences in height.\nAccording to the another aspect of the present invention, focus is attained in accordance with the shape of the recessed part by an operation of the field curvature characteristics in the Under direction, or in accordance with the shape of the protruded part by the operation in the Over direction. An amount of field curvature when the protruded part or the recessed part is in focus can be grasped from a state of movement of the movable lens, and information on differences in height of the central part and the peripheral part of the protruded part or the recessed part can be operated from this amount of field curvature. This information on differences in height is caused to display on the monitor in, for example, the sectional shape and height difference value, or contour and height difference value by output means or the display control means and can be outputted to equipment such as a printer.\nEven in this case, there is provided judging means for dividing an imaging area for the objective lens group into a plurality and judging whether or not the central part and the peripheral part are in focus within a predetermined range by comparing information from these divided areas, and when the central part and the peripheral part are in focus within a predetermined range on the basis of judgment by this judging means, the information on differences in height is preferably displayed, whereby it is possible to maintain accurate information on differences in height."}
{"text": "Glaucoma, which some estimate affects 2 million adults over 40, is an impairment of vision caused by too much fluid pressure within the eye.\nSurgical treatment for glaucoma does not represent a long term cure. While it is effective, it is also expensive and traumatic, and some surgeons will use surgery only as a last resort.\nCarbonic anhydrase inhibitors, prescribed orally work well to treat this disease, but they carry a host of side effects, from nausea to kidney stones. The preferred method of treatment for the disease is instillation by drops to the eye; however, carbonic anhdrase inhibitors have not proven effective when given this way.\nGlaucoma stems from an excess of fluid behind the cornea, the three-layered tissue that acts as a window to let light enter. Fluid carrying nutrients such as potassium and glucose constantly wash the inside of the cornea to keep it healthy, much as tears wash the outside of the cornea.\nIn some middle-aged adults, fluids build up faster than can be absorbed back into the blood, for one of two reasons: the ciliary body (a tiny tissue behind the iris) may excrete too much fluid, or the fluid may not drain off at the normal rate.\nEither way, the excess fluid damages the optic nerve. At first, a glaucoma victim usually experiences a subtle loss of peripheral vision--objects will seem to disappear from certain spots to the side. But glaucoma often leads to middle-age blindness.\nUnfortunately, the two approaches to general drug usage in treating glaucoma--topical (dropped into the eye) and oral--each have a peculiar set of side effects.\nTo make the long journey, oral drugs must be dosed in very high concentration. One class of drugs, called carbonic anhydrase inhibitors when taken orally slow the formation of fluid by inhibiting a chemical reaction at the ciliary body. Along with their well-tested effectiveness, comes nausea, tingling in fingers and toes and other side effects. Oral drugs generally do not, however, cause side effects in the eye.\nCertain topical drugs, other than carbonic anhaydrase inhibitors, e.g., pilocarpine, while causing less systemic effects, can cause severe headaches and constrict the pupil, making the daytime appear dark.\nAccordingly, there is a real and continuing need to develop an inhibitor drug that can be dropped into the eye instead of shallowed, thereby avoiding the present side effects.\nIt is a primary objective of the present invention to develop a highly effective topical carbonic anhydrase inhibitor drug (which previously is only effective orally) for treatment of glaucoma to reduce intraocular eye pressure, and at the same time, avoid the systemic side effects, commonly caused by oral drugs.\nAnother objective of the present invention is to develop a drug for topical treatment of glaucoma, which is not only effective, but which will also pass through the three layered cornea and still be effective enough to work on the ciliary body.\nAnother objective of the present invention is to develop a highly effective, topical drug treatment for glaucoma which is substantially non-harmful to the eye when topically applied.\nAn even further objective of the present invention is to develop an eye treating topical composition which is effective for glaucoma treatment.\nA still further objective is to provide a convenient method of synthesis of certain new and novel compounds which are highly effective topical treatments for glaucoma.\nA further specific objective of the present invention is to provide as a novel compound, N-methyl-2-acetylamino-1,3,4-thiadiazole-5-sulfonamide (methyl acetazolamide), which in pharmaceutically effective amounts is a highly effective topical composition for eye drop treatment of glaucoma.\nThe method and manner of achieving each of the above objectives, as well as others, will become apparent from the detailed description of the invention which follows hereinafter."}
{"text": "The prior art contains many teachings pertaining to toggle type closure fasteners. Some examples of prior United States patents of general interest in relation to this invention are the follwing, made of record herein under 37 C.F.R. 1.56:\nU.S. Pat. No. 497,445 PA1 U.S. Pat. No. 567,621 PA1 U.S. Pat. No. 783,338 PA1 U.S. Pat. No. 1,339,174 PA1 U.S. Pat. No. 1,863,863 PA1 U.S. Pat. No. 1,899,822 PA1 U.S. Pat. No. 3,109,675 PA1 U.S. Pat. No. 3,145,038 PA1 U.S. Pat. No. 3,162,419 PA1 U.S. Pat. No. 3,174,784 PA1 U.S. Pat. No. 3,534,992 PA1 U.S. Pat. No. 3,706,467 PA1 U.S. Pat. No. 4,049,301\nA prime object of the invention is to improve on the known prior art through provision of a closure fastener which is non-handed in comparison to conventional casement and awning type window locks. The mechanism embodied in the invention includes a shallow lock case and filler element which can be flush mounted on one side of a casement window frame and the like so that no parts of the mechanism project into the window sight line when the lock is closed and no parts project into the room to cause interference with window hangings in any position of the mechanism. The entire mechanism is simple, extremely strong and durable and very compact.\nAnother important object of the invention is to provide an action which is very powerful for drawing a sash or other closure inwardly or forcing it outwardly to release a stuck sash. This action involves an over-dead-center mechanism which inhibits forced entry because of a tendency of the mechanism to tighten the locking action rather than release it in response to opening pressure on the closure or sash.\nAnother object and feature of the invention is the provision in the fastener or lock of an essentially straight line draw-in motion which is not critical in relation to alignment with the keeper or strike mounted on the sash.\nOther objects and advantages of the invention will appear to those skilled in the art during the course of the following detailed description."}
{"text": "There are many situations in which it is desirable to transmit, via documents, highly confidential information, such as a PIN through the mail. When such information is printed, it is of course necessary that the information not be visible except when exposed by the user, and there must be some mechanism for identifying to the user if the information has been tampered with so that the PIN can be cancelled and a new PIN issued.\nAccording to the present invention, a tamperproof label assembly for camouflaging or clearly evidently displaying confidential indicia is provided, as well as a method of producing label assemblies obscuring confidential information that is part of each label assembly, while allowing ready exposure of the confidential information in a tamperproof manner. The invention achieves its objectives by using conventional pseudo adhesion layers as part of the label so that once the label components have been separated, to expose the PIN, it is virtually impossible to lay the elements back together in the manner that they were originally provided. The desirable objectives are also facilitated, however, by forming slits and/or perforations in the face layer of the label assembly which cause the label assembly to come apart when the PIN is exposed, making it even more difficult to properly lay the various elements back together again without exposing the tampering, and this is even more particularly facilitated when readable indicia is printed on the top surface of the face ply and at least one slit or perforation extends through the readable indicia.\nThe invention also ensures confidentiality by providing security pantographs on the label face plies so that the underlying PIN cannot be read, the security pantographs typically provided by first printing in black or blue ink, and then printing with the other of those colors. Also, the slits form a triangularly shaped element in the face ply of a distinctly different color (e.g. red) from the security pantographs which allows ready access to the components for separating various plies and layers at the pseudo adhesion layers.\nTo facilitate easy production of the label assemblies according to the invention, each label assembly typically includes an image ply on which the PIN is generated. The image is formed on the impact ply by applying an uninked impact printing element (stylus) to the top surface of the face ply, while the security pantographs are applied to the face ply by non-impact printing (e.g. printing plates), so as not to form images on the underlying impact ply.\nAccording to one aspect of the present invention a method of producing label assemblies is provided comprising the following steps: (a) Producing a web of label components having a face ply with a top surface, underlying pseudo adhesion layers, an underlying image ply, an underlying adhesive layer, and a bottom release paper carrier ply. (b) Acting on the web face ply to non-impact print a security pantograph thereon. (c) Acting on the web face ply with an impact element to generate a confidential information image on the image ply, which image is not visible from the face ply. And, (d) die cutting labels out of the web, and removing matrix material formed by die cutting so that the web then comprises the bottom release paper carrier ply with a plurality of spaced labels thereon, each having a confidential information image therein.\nThe method also preferably comprises the further steps of imaging readable indicia on the face ply (not overlying the confidential information image), and forming a plurality of slits and/or perforations in the face ply and one pseudo adhesion layer, including through the readable indicia, to facilitate the tamperproof functionality of the labels. The first security pantograph is typically printed in black or blue ink, and then the security pantograph is printed in the other of blue or black ink. A diagonal slit forms a generally triangularly shaped element in the face ply which preferably is printed with a distinctly different color, such as red, and which is readily (and easily discernible) removed to allow access to the face ply to separate the pseudo adhesion layers.\nAccording to another aspect of the present invention a tamperproof label assembly is provided comprising the following elements: An image ply having top and bottom surfaces, the top surface for generating an image in the form of confidential indicia when impacted with an impact printer stylus. Permanent adhesive operatively associated with the bottom surface of the image ply. A face ply having top and bottom surfaces, and having a security pantograph imaged on the top surface thereof. And, pseudo adhesion layers disposed between the face ply bottom surface and the image ply top surface for releasably holding the face and image plies together, but when separated exposing the image ply top surface without damage thereto, and once separated not capable of re-adhesion in as effective a manner as prior to separation.\nConfidential indicia is imaged on the top surface of the image ply. The assembly also typically comprises security slits and/or perforations formed on the face ply and one of the pseudo adhesion layers for effecting separation of the portions of the second face ply when pseudo adhesion layers are separated, precluding effective reattachment of the pseudo adhesion layers when separated due to that separation. The security slits may include a diagonal slit cooperating with a longitudinal slit and defining a triangularly shaped portion of the face ply, with the top surface of the triangularly shaped portion having printing thereon of distinctly different color (e.g. red) than the security pantograph. Readable indicia is typically imaged on the top surface of the face ply and at least one of the slits or perforations passes through the readable indicia causing the face ply to separate into different portions each containing only part of the readable indicia, when the pseudo adhesion layers are separated. The security slits and perforations may include a plurality of generally L-shaped slits, at least one longitudinal slit, and at least one longitudinal perforation.\nTypically the label assembly also includes a liner ply between the bottom surface of the image and the permanent adhesive; a release sheet engaging the permanent adhesive; a permanent adhesive between the top surface of the image ply and one of the pseudo adhesion layers; and a permanent adhesive between the bottom surface of the face ply and one of the pseudo adhesion layers.\nThe invention also comprises a tamperproof label assembly comprising the following elements: An image ply having top and bottom surfaces, the top surface for generating an image in the form of confidential indicia when impacted with an impact printer stylus. A face ply having top and bottom surfaces, and having a security pantograph imaged on the top surface thereof. Pseudo adhesion layers disposed between the face ply bottom surface and the image ply top surface for releasably holding the face and image plies together, but when separated exposing the image ply top surface without damage thereto, and once separated not capable of re-adhesion in as effective a manner as prior to separation. And, security lines of weakness or parting formed in the face ply and one of the pseudo adhesion layers for effecting separation of portions of the face ply when the pseudo adhesion layers are separated, precluding effective reattachment of the pseudo adhesion layers once separated due to the separation.\nIt is the primary object of the present invention to provide for the effective transmission of confidential information, yet allowing ready exposure of the confidential information but in a tamperproof manner, i.e. so that if the confidential information is once exposed that is thereafter ever apparent. This and other objects of the invention will become clear from an inspection of the detailed description of the invention, and from the appended claims."}
{"text": "1. Field of the Invention\nThis invention relates to an exerciser, more particularly to a hand-held exerciser.\n2. Description of the Related Art\nReferring to FIG. 1, a conventional exerciser is shown to comprise a fixed body 1, a movable body 2, a coiled spring 3 and a spring adjustment device 4. The fixed body 1 has an elongated base plate (1a), an upright outer tube (1b) connected to the base plate (1a) at its lower end, and a fastening device (1c) which is mounted under the lower end of the outer tube (1b) in order to connect rotatably the outer tube (1b) to the fixed body 1. The lower side of the base plate (1a) defines substantially a planar, flat support surface in order to be disposed stably on the ground. The movable body 2 has a transverse press bar (2a) and an inner tube (2b) connected perpendicularly to the transverse bar (2a) at its upper end. The lower portion of the inner tube (2a) is received slidably in the outer tube (1b). The coiled spring 3 is disposed in the inner tube (2a). The spring adjustment device 4 has a threaded rod (4c) which has a head (4d) that engages a rotary knob (4a), and a nut member (4b) connected to its lower end. The spring force of the coiled spring 3 can be increased or decreased by means of rotating the rotary knob (4a) and thereby the threaded shaft (4c), causing the nut member (4b) to move relative to the threaded shaft (4c) and to compress or release the coiled spring 3.\nIn use, the user may place the base plate (1a) of the fixed body 1 on the ground and depress the movable body 2 downward against the spring force of the coiled spring 3 for exercising purposes. However, because the weight of the user's body facilitates the depression of the movable body 2 when the user normally bends his/her body to depress the movable body 2 downward by means of arms, the user cannot simply exercise his/her arms by means of the conventional exerciser."}
{"text": "With a diesel engine driving an automobile, it is usually necessary to heat up glow plugs to a high temperature of, for example, eight hundred degrees centigrade which is suitable for quickly starting the diesel engine, prior to performing the starting operation of the engine by means of a known starter motor. However, under certain environmental temperature conditions of a diesel engine, for example, under a condition where the diesel engine per se is sufficiently warmed up, it is possible to start the heating of the glow plugs at the same time as operating the starter motor which operation is performed by shifting a key switch to a starting position thereof.\nA conventional control apparatus for glow plugs disclosed in, for example, U.S. Pat. No. 3,675,033 belonging to Richard et al does not always take such environmental temperature condition of a diesel engine into consideration for achieving control of the heating of glow plugs."}
{"text": "The present invention relates to a certain heterocyclic pyrimidine and its use in processes and compositions for altering the flavors and aromas of various materials such as tobaccos, foodstuffs, and the like, as well as the novel pyrimidines and processes for producing them.\nBecause of the tremendous consumption of foods, tobaccos, and other materials, there has been an increasing interest in substances and methods for imparting flavors to such consumable materials. This interest has been stimulated not only because of the inadequate quantity of natural flavoring materials available, but perhaps even more importantly, because of the need for materials which can combine several nuances, will be more stable than natural materials, will blend better with other flavors or flavoring composition components, and will generally provide superior products.\nVarious pyrimidines containing 4-substituents have been shown in the art, but no flavoring, enhancing or other organoleptic properties have been shown or suggested. Thus, U.S. Pat. No. 3,272,811 shows a variety of pyrimidine."}
{"text": "Archery has become increasingly popular over the past few years as evidenced by various tournaments and competitions where archers shoot at targets either in an indoor or outdoor environment in order to determine their level of skill. Additionally, game hunting with bows and arrows is on the rise and, usually, the archer is in an elevated position in a tree stand where he can survey the ground area around him. As the hunter has few opportunities to shoot at game, it is important that he determine the distance to the game in a quick and efficient manner. This distance determination is important when a compound bow is being used as this type of bow is equipped with a string peep sight and sight pins which have been pre-set at known yardages, such as 20, 30 etc. yards in order for the archer to easily determine the distance to the game and thereby assist him in the aiming and shooting of the arrow. Therefore, any assistance that can be provided to determine which sight pin should be selected would increase his chance of success.\nHeretofore, these distances were determined by the hunter in various ways. He could, at some point beforehand, pace off the radial distances from the tree in which the stand is located and then use landmarks or place markers to establish/designate distances therefrom. He could also rely on commercially available optical rangefinders, which are both cumbersome and costly, or he could just guess at the distance.\nThe present invention is directed to a simple hand held device which can be suspended from the hunter's hand, and by eyeing fixed and adjustable sight pins he can readily determine the distance to the target and select the proper bow sight pin accordingly."}
{"text": "The present disclosure is related to reel based closure devices for various articles, such as braces, medical devices, shoes, clothing, apparel, and the like. Such articles typically include some closure system, which allows the article to be placed about a body part and closed or tightened about the body part. The closure systems are typically used to maintain or secure the article about the body part. For example, shoes are typically placed over an individual's foot and a shoelace is tensioned and tied to close and secure the shoe about the foot. Conventional closure systems have been modified in an effort to increase the fit and/or comfort of the article about the body part. For example, shoe lacing configurations and/or patterns have been modified in an attempt to increase the fit and/or comfort of wearing shoes. Conventional closure systems have also been modified in an effort to decrease the time in which an article may be closed and secured about the body part. These modifications have resulted in the use of various pull cords, straps, and tensioning devices that enable the article to be quickly closed and secured to the foot."}
{"text": "Current color selection technologies mainly include two methods. Method 1, as illustrated in FIG. 1, is selecting a color by entering a value representing the color, i.e., a hexadecimal value corresponding to a legal color name. This method is not useful because most people do not know each color's value. If a user does not enter the precise value corresponding to the desired color, the obtained color can be different. Method 2, as illustrated in FIG. 2, is selecting a color from a color palette. This method can usually be found in mobile phone applications. However, due to limited screen sizes of mobile devices, a larger number of color options, and a large finger press area size, it can be difficult to accurately select a color of desire without multiple attempts. Thus, the second method can be inefficient and inaccurate due to its complicated operation requirement."}
{"text": "A virtual machine (VM) is a software implementation of a physical computer. Computer programs designed to execute on the physical machine execute in a similar way when executed on a VM. A VM provides a complete system platform to support a full operating system (OS). A physical machine can be shared between users by using different VMs, each running a different OS.\nModern processor architectures have enabled virtualization techniques that allow multiple operating systems and VMs to run on a single physical machine. These techniques use a hypervisor layer that runs directly on the physical hardware and mediates accesses to physical hardware by providing a virtual hardware layer to the operating systems running in each virtual machine. The hypervisor can operate on the physical machine in conjunction with a ‘native VM’. Alternatively, the hypervisor can operate within an operating system running on the physical machine, in conjunction with a ‘hosted VM’ operating at a higher software level.\nExamples of VM technology are:                Linux Kernel-Based Virtual Machine (KVM) allows one or more Linux or Windows virtual machines to be run on top of an underlying Linux that runs KVM.        Xen allows a guest (virtualized) Linux to be run on top of Linux.        Parallels allows Linux and Windows on top of Mac OS X.        VMWare allows Linux and Windows systems on top of Mac OS X, Windows and Linux systems.        \nA user may want to migrate a workload operating on one physical machine (host A) to another physical machine (host B), for example, for machine maintenance or for performance optimisation. If the instruction set architecture is the same on both host A and host B, the VM needs to be shut down on host A, restarted on host B, and the workload migrated. However, if the ISA on each physical machine is different, a migration is problematic, because, for example, the format state of the VM on host A is inappropriate for the format state of the VM on host B."}
{"text": "This invention relates to a method for controlling the pressure output of an engine-driven pump system. Specifically, this invention relates to a method of controlling the discharge pressure of an engine-driven pump for use in a fire truck.\nIt is vital to control the discharge pressure of an engine-driven fire pump mounted on or in a fire truck. The pump must supply water at various rates and steady pressure so that firemen operating the hoses at a fire scene can control the reaction force generated by their hose nozzles. Fire pumps as used here are centrifugal pumps. These pumps add pressure to the incoming source of water. Therefore pressure changes in the supply are pressure changes in the discharge. This is problematic because even slight variations in pressure in the supply line leading to the intake of the pump are amplified by the pump on the discharge side, causing surges or oscillations in the water flow discharge at the nozzle and corresponding changes in the reaction forces. Such changes are extremely dangerous, as they can pull a nozzle out of the fireman's grip, or even throw him or her off a ladder or ledge.\nThe simplest prior art device for controlling the pressure output of the fire pump is a mechanical relief valve which opens to discharge excess water when the pressure is higher than the desired output pressure. A shortcoming of such a valve, however, is that, the relief valve only functions to dissipate excess pressure, and has no utility in situations where the pressure is too low, such as when the water source is being depleted or another hose is connected to the system. In addition, if the pump engine continues to operate at full speed after the relief valve is opened, water will be continuously recirculated in the system, resulting in needless waste and wear and tear on the pump and engine. Overheating of the pump and engine is also more likely.\nElectronically operated pressure controlled systems have been developed. Two such systems are disclosed in U.S. Pat. Nos. 3,786,869 and 4,189,005 to McLoughlin, the subject matter of which is herein incorporated by reference. In these systems, the desired output pressure is dialed in or otherwise transmitted to a control box on the board of the fire truck, where it is compared to the actual output pressure as measured by a transducer. Any difference between the desired and actual output pressure is converted to an electrical signal which is fed to a DC motor which increases or decreases the rpm of the centrifugal pump as needed until the desired output pressure is reached. A shortcoming of this type of system is that, because the response time of the servo-mechanism controlling the engine is slow, much time can pass before the appropriate rpm and correct discharge pressure are reached. This is especially troublesome during transient events, such as overpressure spikes, where the system's response time is greater than the length of the event. Furthermore, no allowance is made for situations such as when the engine is already at idle and the incoming pressure suddenly increases, or is higher than desired, such as what can happen when the pump is connected to a hydrant. Recent engine technology has replaced the servo with direct commands to the engine computer or an electrical throttle control which can improve response times.\nAnother control system of interest is disclosed in German Patent No. 1,274,402 to Mueller and Company, which discloses an engine-driven pump which responds to an over pressure in the supply line by simultaneously opening a pressure relief valve and mechanically reducing the engine speed. The shortcoming of this purely mechanical system is that by its nature, in cases of over pressure, the relief valve will always be open to some extent, allowing some fluid to always bypass the relief valve, and the engine rpm will always be above its idle setting to a certain extent.\nAnother pressure control system of interest is disclosed in U.S. Pat. No. 5,888,051, which discloses an engine-driven pump which responds to an over pressure in the supply line and lowers the engine rpm and simultaneously controls a pressure relief valve which may be commanded to open and dump water for short durations to relieve over pressure spikes, or for longer duration to relieve excess water coming into the pump. The shortcomings of this system are that the change in engine speed and the relief valve may be operated at the same time resulting in a waste of water. In addition, operating the engine and relief valve simultaneously results in a needlessly complicated response system.\nAccordingly, a need exists for a new and improved electronically operated fire pump discharge pressure control system for quickly and safely responding to drops or increases in the incoming pressure of a fire pump, which change the discharge pressure required, as well as changes in discharge pressure due to the opening or shutting off of various valves downstream of the pump."}
{"text": "1. Field\nExample embodiments relate to methods of forming contact holes.\n2. Description of the Related Art\nPhotolithography has been used for forming a pattern of a semiconductor device. However, it is difficult to form a pattern having a minute pitch below 40 nm due to limits in the resolution. Thus, a double patterning method was developed. However, the double patterning method may have too many steps and/or may be complicated.\nThus, a method of forming a minute pitch pattern, particularly, contact holes using a direct self assembly (DSA) has been developed. However, in the DSA method, a material for the DSA may be arranged in a hexagonal shape for a thermodynamic stability such that it may be difficult to form contact holes having a desired arrangement."}
{"text": "Serial Attached SCSI (SAS) is a term referring to various technologies designed to implement data transfer between computer devices. The SAS protocol is a serial successor to the parallel Small Computer System Interface. In the SAS protocol, all SAS devices are either an initiator device, a target device, or an expander device. Initiator devices are devices that begin an SAS data transfer, while target devices are the devices to which initiator devices transfer data. Together, initiator devices and target devices are known as end devices.\nSAS expanders are devices that facilitate data transfer between multiple initiator devices and multiple target devices. The SAS protocol utilizes a point-to-point bus topology. Therefore, if an initiator device is required to connect to multiple target devices, a direct connection must be made between the initiator device and each individual target device in order to facilitate each individual data transfer between the initiator device and each individual target device. SAS expanders manage the connections and data transfer between multiple initiator devices and multiple target devices. SAS expanders may contain SAS devices."}
{"text": "The present invention relates to a spin exciting method, a magnetic resonance imaging method, and a magnetic resonance imaging system. More particularly, the present invention relates to a spin exciting method, a magnetic resonance imaging method, and a magnetic resonance imaging system for performing magnetic resonance imaging according to a fast spin echo (FSE) technique combined with an inversion recovery (IR) technique.\nIn a magnetic resonance imaging (MRI) system, a subject of imaging is carried into a bore of a magnet system, that is, an imaging space in which a static magnetic field is created. Magnetic field gradients and a radio-frequency magnetic field are applied to the subject in order to excite spins in the subject. Consequently, a magnetic resonance signal is induced, and an image is reconstructed based on the signal received.\nA sequence of exciting spins so as to induce a magnetic resonance signal and receiving the signal is repeated at predetermined intervals of a repetition time TR. The TR is often set to a time long enough for the excited spins to recover to exhibit an original longitudinal magnetization. When an imaging time must be shortened, the TR is set to a short time and spins are forcibly recovered. The forcible recovery of spins is achieved with additional excitation. This technique is referred to as fast recovery.\nJapanese Examined Patent Publication No. 4-21488 describes that the fast recovery is combined with the IR. In short, as shown in FIG. 7, spins are turned 180xc2x0 with application of a 180xc2x0 pulse and thus brought to an inversion. Thereafter, when a predetermined inversion time TI has elapsed, a 90xc2x0 pulse is applied in order to turn the spins 90xc2x0. A free induction decay (FID) signal that is induced accordingly is then acquired.\nThereafter, when a half of an echo time TE has elapsed, a xe2x88x92180xc2x0 pulse is applied in order to reverse the spins. Thereafter, when a half of the TE has elapsed, a xe2x88x9290xc2x0 pulse is applied in order to turn the spins xe2x88x9290xc2x0 and a 180xc2x0 pulse is then applied in order to reverse the spins. Thus, the fast recovery of the spins is achieved.\nU.S. Pat. No. 6,054,853 describes that the fast recovery is combined with the FSE. In short, as shown in FIG. 8, a 90xc2x0 x pulse is applied in order to excite spins and turn them 90xc2x0 with respect to an x axis. Thereafter, when a half of an echo space esp has elapsed, a 180xc2x0 y pulse is applied in order to reverse the spins with respect to a y axis. Thereafter, when the esp has elapsed, the 180xc2x0 y pulse is applied in order to reverse the spins again with respect to the y axis. When the echo space esp has elapsed, the 180xc2x0 y pulse is applied in order to reverse the spins again with respect to the y axis. Consequently, a spin echo is acquired during the echo space esp between applications of the 180xc2x0 y pulse.\nWhen a half of the echo space esp has elapsed since the last application of the 180xc2x0 y pulse, a xe2x88x9290xc2x0 x pulse is applied in order to turn the spins xe2x88x9290xc2x0, and a 180xc2x0 x pulse is applied in order to reverse the spins. Thus, the fast recovery of the spins is achieved.\nAccording to the technology described in the Japanese Examined Patent Publication No. 4-21488, the fast recovery employing the xe2x88x9290xc2x0 pulse and 180xc2x0 pulse is achieved through unselective excitation that is not intended to select a slice. This disables multiple slice imaging that interleaves pulse sequences like the foregoing one and involves a plurality of slices.\nAccording to the related art described in the U.S. Pat. No. 6054853, the fast recovery employing the xe2x88x9290xc2x0 x pulse and 180xc2x0 is performed through selective excitation. However, a selected slice is not perfectly square. It is not easy to properly achieve the fast recovery using the two selective excitation pulses.\nMoreover, the number of reversals of spins stemming from application of the 180xc2x0 y pulse is an odd value. If the degree of the reversal of spins stemming from application of the 180xc2x0 y pulse has an error, the spins are not restored to exactly face along the x-y plane. Therefore, the succeeding fast recovery is achieved imperfectly.\nTherefore, an object of the present invention is to realize a spin exciting method, a magnetic resonance imaging method, and a magnetic resonance imaging system for properly performing fast recovery during magnetic resonance imaging in which the fast spin echo technique combined with the inversion recovery technique is implemented.\n(1) In one aspect of the present invention intended to solve the aforesaid problems, there is provided a spin exciting method for producing an image using a magnetic resonance signal induced by spins in a subject being imaged according to the fast spin echo technique combined with the inversion recovery. Specifically, a 180xc2x0 pulse is applied in order to excite spins. Thereafter, when a first time has elapsed, a first 90xc2x0 x pulse is applied in order to excite the spins. Thereafter, when a second time has elapsed, a 180xc2x0 y pulse is applied in order to excite the spins. Thereafter, when a third time that is double the second time has elapsed, the 180xc2x0 y pulse is applied an odd number of times in order to sequentially excite the spins. Thereafter, when the second time has elapsed, a second 90xc2x0 x pulse is applied in order to excite the spins.\n(2) In another aspect of the present invention intended to solve the aforesaid problems, there is provided a magnetic resonance imaging method for producing an image using a magnetic resonance signal induced by spins in a subject being imaged according to the fast spin echo technique combined with the inversion recovery technique. Specifically, a 180xc2x0 pulse is applied in order to excite spins. Thereafter, when a first time has elapsed, a first 90xc2x0 x pulse is applied in order to excite the spins. Thereafter, when a second time has elapsed, a 180xc2x0 y pulse is applied in order to excite the spins. Thereafter, when a third time that is double the second time has elapsed, the 180xc2x0 y pulse is applied an odd number of times in order to sequentially excite the spins. Thereafter, when the second time has elapsed, a second 90xc2x0 x pulse is applied in order to excite the spins. A spin echo is read during the third time, and an image is produced based on the spin echo.\n(3) In another aspect of the present invention intended to solve the aforesaid problems, there is provided a magnetic resonance imaging system for producing an image using a magnetic resonance signal induced by spins in a subject being imaged according to the fast spin echo technique combined with the inversion recovery technique. The magnetic resonance imaging system includes a spin exciting means, an echo reading means, and an image producing means. The spin exciting means excites spins with application of a 180xc2x0 pulse. Thereafter, when a first time has elapsed, the spin exciting means applies a first 90xc2x0 x pulse to excite the spins. Thereafter, when a second time has elapsed, the spin exciting means applies a 80xc2x0 y pulse to excite the spins. Thereafter, when a third time that is double the second time has elapsed, the spin exciting means applies the 180xc2x0 y pulse an odd number of times to excite the spins sequenlially. Thereafter, when the second time has elapsed, the spin exciting means applies a second 90xc2x0 x pulse to excite the spins. The echo reading means reads a spin echo during the third time. The image producing means produces an image according to the spin echo.\nIn the aspects of the present invention set forth in paragraphs (1) to (3), a 180xc2x0 pulse is applied in order to excite spins. Thereafter, when the first time has elapsed, a first 90xc2x0 x pulse is applied in order to excite the spins. Thereafter, when the second time has elapsed, a 180xc2x0 y pulse is applied in order to excite the spins. Thereafter, when the third time that is double the second time has elapsed, the 180xc2x0 y pulse is applied an odd number of times in order to excite the spins. Thereafter, when the second time has elapsed, a second 90xc2x0 x pulse is applied in order to excite the spins. The spins are reversed an even number of times due to the applications of the 180xc2x0 y pulse, and therefore returned to exactly face along the xy plane. Consequently, the spins are accurately recovered with the subsequent application of the 90xc2x0 x pulse.\nMoreover, the inversion recovery is performed with application of the 90xc2x0 x pulse alone. Therefore, imperfect fast recovery stemming from the employment of two selective excitation pulses as conventionally will not take place.\nIn order to successfully achieve fast recovery of spins whose relaxation time is relatively long, the second 90xc2x0 x pulse should preferably be a +90xc2x0 x pulse.\nIn order to achieve imaging while enhancing a magnetic resonance signal induced by spins whose relaxation time is relatively long, the first time should preferably be shorter than the polarity recovery time for spins that are used for imaging.\nIn order to successfully achieve fast recovery of spins whose relaxation time is relatively short, the second 90xc2x0 x pulse should preferably be a xe2x88x9290xc2x0 x pulse.\nIn order to achieve imaging while enhancing a magnetic resonance signal induced by spins whose relaxation time is relatively short, the first time should preferably be longer than the polarity recovery time for spins that are used for imaging.\nIn order to produce a tomographic image, excitation should preferably be selective excitation all the time.\nIn order to achieve multiple slice imaging, a series of excitations starting with excitation initiated with application of the 180xc2x0 pulse and ending with excitation initiated with application of the second 90xc2x0 x pulse should preferably be started at successive time instants within the first time with slices changed sequentially.\nAccording to the present invention, there are provided a spin exciting method, a magnetic resonance imaging method, and a magnetic resonance imaging system for making it possible to properly achieve fast recovery during magnetic resonance imaging in which the fast spin echo technique combined with the inversion recovery technique is implemented.\nFurther objects and advantages of the present invention will be apparent from the following description of the preferred embodiments of the invention as illustrated in the accompanying drawings."}
{"text": "The invention pertains to the field of computer directed instruments for performing the polymerase chain reaction (hereinafter PCR). More particularly, the invention pertains to automated instruments that can perform the polymerase chain reaction simultaneously on many samples with a very high degree of precision as to results obtained for each sample. This high precision provides the capability, among other things, of performing so-called \"quantitative PCR\".\nTo amplify DNA (Deoxyribose Nucleic Acid) using the PCR process, it is necessary to cycle a specially constituted liquid reaction mixture through a PCR protocol including several different temperature incubation periods. The reaction mixture is comprised of various components such as the DNA to be amplified and at least two primers selected in a predetermined way to as to be sufficiently complementary to the sample DNA as to be able to create extension products of the DNA to be amplified. The reaction mixture includes various enzymes and/or other reagents, as well as several deoxyribonucleoside triphosphates such as dATP, dCTP, dGTP and dTTP. Generally, the primers are oligonucleotides which are capable of acting as a point of initiation of synthesis when placed under conditions in which synthesis of a primer extension product which is complimentary to a nucleic acid strand is induced, i.e., in the presence of nucleotides and inducing agents such as thermostable DNA polymerase at a suitable temperature and pH.\nThe Polymerase Chain Reaction (PCR) has proven a phenomenally successful technology for genetic analysis, largely because it is so simple and requires relatively low cost instrumentation. A key to PCR is the concept of thermocycling: alternating steps of melting DNA, annealing short primers to the resulting single strands, and extending those primers to make new copies of double stranded DNA. In thermocycling, the PCR reaction mixture is repeatedly cycled from high temperatures (&gt;90.degree. C.) for melting the DNA, to lower temperatures (40.degree. C. to 70.degree. C.) for primer annealing and extension. The first commercial system for performing the thermal cycling required in the polymerase chain reaction, the Perkin-Elmer Cetus DNA Thermal Cycler, was introduced in 1987.\nApplications of PCR technology are now moving from basic research to applications in which large numbers of similar amplifications are routinely run. These areas include diagnostic research, biopharmaceutical development, genetic analysis, and environmental testing. Users in these areas would benefit from a high performance PCR system that would provide the user with high throughput, rapid turn-around time, and reproducible results. Users in these areas must be assured of reproducibility from sample-to-sample, run-to-run, lab-to-lab, and instrument-to-instrument.\nFor example, the physical mapping process in the Human Genome Project may become greatly simplified by utilizing sequence tagged sites. An STS is a short, unique sequence easily amplified by PCR and which identifies a location on the chromosome. Checking for such sites to make genome maps requires amplifying large numbers of samples in a short time with protocols which can be reproducibly run throughout the world.\nAs the number of PCR samples increases, it becomes more important to integrate amplification with sample preparation and post-amplification analysis. The sample vessels must not only allow rapid thermal cycling but also permit more automated handling for operations such as solvent extractions and centrifugation. The vessels should work consistently at low volumes, to reduce reagent costs.\nGenerally PCR temperature cycling involves at least two incubations at different temperatures. One of these incubations is for primer hybridization and a catalyzed primer extension reaction. The other incubation is for denaturation, i.e., separation of the double stranded extension products into single strand templates for use in the next hybridization and extension incubation interval. The details of the polymerase chain reaction, the temperature cycling and reaction conditions necessary for PCR as well as the various reagents and enzymes necessary to perform the reaction are described in U.S. Pat. Nos. 4,683,202, 4,683,195, EPO Publication 258,017 and 4,889,818 (Taq polymerase enzyme patent), which are hereby incorporated by reference.\nThe purpose of a polymerase chain reaction is to manufacture a large volume of DNA which is identical to an initially supplied small volume of \"seed\" DNA. The reaction involves copying the strands of the DNA and then using the copies to generate other copies in subsequent cycles. Under ideal conditions, each cycle will double the amount of DNA present thereby resulting in a geometric progression in the volume of copies of the \"target\" or \"seed\" DNA strands present in the reaction mixture.\nA typical PCR temperature cycle requires that the reaction mixture be held accurately at each incubation temperature for a prescribed time and that the identical cycle or a similar cycle be repeated many times. A typical PCR program starts at a sample temperature of 94.degree. C. held for 30 seconds to denature the reaction mixture. Then, the temperature of the reaction mixture is lowered to 37.degree. C. and held for one minute to permit primer hybridization. Next, the temperature of the reaction mixture is raised to a temperature in the range from 50.degree. C. to 72.degree. C. where it is held for two minutes to promote the synthesis of extension products. This completes one cycle. The next PCR cycle then starts by raising the temperature of the reaction mixture to 94.degree. C. again for strand separation of the extension products formed in the previous cycle (denaturation). Typically, the cycle is repeated 25 to 30 times.\nGenerally, it is desirable to change the sample temperature to the next temperature in the cycle as rapidly as possible for several reasons. First, the chemical reaction has an optimum temperature for each of its stages. Thus, less time spent at nonoptimum temperatures means a better chemical result is achieved. Another reason is that a minimum time for holding the reaction mixture at each incubation temperature is required after each said incubation temperature is reached. These minimum incubation times establish the \"floor\" or minimum time it takes to complete a cycle. Any time transitioning between sample incubation temperatures is time which is added to this minimum cycle time. Since the number of cycles is fairly large, this additional time unnecessarily lengthens the total time needed to complete the amplification.\nIn some prior automated PCR instruments, the reaction mixture was stored in a disposable plastic tube which is closed with a cap. A typical sample volume for such tubes was approximately 100 microliters. Typically, such instruments used many such tubes filled with sample DNA and reaction mixture inserted into holes called sample wells in a metal block. To perform the PCR process, the temperature of the metal block was controlled according to prescribed temperatures and times specified by the user in a PCR protocol file. A computer and associated electronics then controlled the temperature of the metal block in accordance with the user supplied data in the PCR protocol file defining the times, temperatures and number of cycles, etc. As the metal block changed temperature, the samples in the various tubes followed with similar changes in temperature. However, in these prior art instruments not all samples experienced exactly the same temperature cycle. In these prior art PCR instruments, errors in sample temperature were generated by nonuniformity of temperature from place to place within the metal sample block, i.e., temperature gradients existed within the metal of the block thereby causing some samples to have different temperatures than other samples at particular times in the cycle. Further, there were delays in transferring heat from the sample block to the sample, but the delays were not the same for all samples. To perform the PCR process successfully and efficiently, and to enable so called \"quantitative\" PCR, these time delays and temperature errors must be minimized to a great extent.\nThe problems of minimizing time delays for heat transfer to and from the sample liquid and minimizing temperature errors due to temperature gradients or nonuniformity in temperature at various points on the metal block become particularly acute when the size of the region containing samples becomes large. It is a highly desirable attribute for a PCR instrument to have a metal block which is large enough to accommodate 96 sample tubes arranged in the format of an industry standard microtiter plate.\nThe microtiter plate is a widely used means for handling, processing and analyzing large numbers of small samples in the biochemistry and biotechnology fields. Typically, a microtiter plate is a tray which is 35/8 inches wide and 5 inches long and contains 96 identical sample wells in an 8 well by 12 well rectangular array on 9 millimeter centers. Although microtiter plates are available in a wide variety of materials, shapes and volumes of the sample wells, which are optimized for many different uses, all microtiter plates have the same overall outside dimensions and the same 8.times.12 array of wells on 9 millimeter centers. A wide variety of equipment is available for automating the handling, processing and analyzing of samples in this standard microtiter plate format.\nGenerally microtiter plates are made of injection molded or vacuum formed plastic and are inexpensive and considered disposable. Disposability is a highly desirable characteristic because of the legal liability arising out of cross contamination and the difficulty of washing and drying microtiter plates after use.\nIt is therefore a highly desirable characteristic for a PCR instrument to be able to perform the PCR reaction on up to 96 samples simultaneously said samples being arranged in a microtiter plate format.\nOf course, the size of the metal block which is necessary to heat and cool 96 samples in an 8.times.12 well array on 9 millimeter centers is fairly large. This large area block creates multiple challenging engineering problems for the design of a PCR instrument which is capable of heating and cooling such a block very rapidly in a temperature range generally from 0 to 100.degree. C. with very little tolerance for temperature variations between samples. These problems arise from several sources. First, the large thermal mass of the block makes it difficult to move the block temperature up and down in the operating range with great rapidity. Second, the need to attach the block to various external devices such as manifolds for supply and withdrawal of cooling liquid, block support attachment points, and associated other peripheral equipment creates the potential for temperature gradients to exist across the block which exceed tolerable limits.\nThere are also numerous other conflicts between the requirements in the design of a thermal cycling system for automated performance of the PCR reaction or other reactions requiring rapid, accurate temperature cycling of a large number of samples. For example, to change the temperature of a metal block rapidly, a large amount of heat must be added to, or removed from the sample block in a short period of time. Heat can be added from electrical resistance heaters or by flowing a heated fluid in contact with the block. Heat can be removed rapidly by flowing a chilled fluid in contact with the block. However, it is seemingly impossible to add or remove large amounts of heat rapidly in a metal block by these means without causing large differences in temperature from place to place in the block thereby forming temperature gradients which can result in nonuniformity of temperature among the samples.\nEven after the process of addition or removal of heat is terminated, temperature gradients can persist for a time roughly proportional to the square of the distance that the heat stored in various points in the block must travel to cooler regions to eliminate the temperature gradient. Thus, as a metal block is made larger to accommodate more samples, the time it takes for temperature gradients existing in the block to decay after a temperature change causes temperature gradients which extend across the largest dimensions of the block can become markedly longer. This makes it increasingly difficult to cycle the temperature of the sample block rapidly while maintaining accurate temperature uniformity among all the samples.\nBecause of the time required for temperature gradients to dissipate, an important need has arisen in the design of a high performance PCR instrument to prevent the creation of temperature gradients that extend over large distances in the block. Another need is to avoid, as much as possible, the requirement for heat to travel across mechanical boundaries between metal parts or other peripheral equipment attached to the block. It is difficult to join metal parts in a way that insures uniformly high thermal conductance everywhere across the joint. Nonuniformities of thermal conductance will generate unwanted temperature gradients."}
{"text": "It is known to provide play sets for use with reduced scale (e.g., 1/64 Hot Wheels® and Matchbox) toy vehicles. Conventional scaled toy vehicle play sets are not well-suited for use by a child while riding in a vehicle such as a car. They are usually too large to set up and use in a car or other vehicle and often lack structure to secure toy vehicles and other loose pieces of the play set during travel. It is believed that a toy vehicle play set adapted for use by a child riding in a vehicle and that can secure its components would be desirable."}
{"text": "1. Field of the Invention\nThe present invention relates generally to processes for ensuring an importer's compliance with domestic customs requirements. More specifically, the invention relates to systems and methods for correcting erroneous information reported to customs agents, as well as the effects of such erroneous reports, after the imported goods to which the reported information refer are received by the importer.\n2. Description of Related Art\nWhen a company imports goods from foreign sources, the company is subject to numerous U.S. Customs regulations. One set of requirements is based upon a complex classification system created by U.S. Customs, known as the Harmonized Tariff Schedule (HTS) code. Within this system, U.S. Customs classifies types of products using 10-digit HTS codes. These codes are used, for example, in determining the applicable tariff rates on different types of products imported into the U.S. An company importing goods must correctly classify imported products under the HTS code. This task is especially difficult for companies that import a large variety of products, such as automobile parts, because of numerous complexities in the classification system.\nLarge importers of foreign goods often employ customs brokers to serve as agents between the importer and U.S. Customs. Brokers fulfill a number of duties related to ensuring the importer's compliance with customs rules, including gathering and delivering information regarding imported goods to customs agents. One such duty is classification of goods under the HTS code. A customs broker working for the importer traditionally classifies imported products by assigning what he decides are the appropriate HTS Codes. The broker also assigns what are known as “attribute classifications,” meaning classifications of other relevant attributes, such as classifications relating to NAFTA Certificates, anti-dumping and Department of Transportation specifications. In addition, a large number of parties may be involved in the importation process. For example, the importer may be working with a large number of foreign manufacturers or suppliers. Each imported item, from every manufacturer and every supplier, must be accounted for within the HTS classification system. In the case of an imported automobile, whose parts may originate from a variety of different overseas sources, vendors and suppliers, providing HTS classification for each part and each component is a daunting and difficult task, with many inherent complexities.\nOnce assigned, the customs broker reports classifications to U.S. Customs and pays duties on behalf of the company based on the reported classifications. Unfortunately, the information reported to U.S. Customs by the customs broker sometimes fails to conform to the actual facts surrounding the goods imported. For example, a broker may not be able to discern precisely how many parts are contained in a carton, or may be unable to accurately classify certain automobile parts that do not clearly fall within one of the Customs-defined HTS codes. For example, an automobile typically includes a transmission system which, in turn, includes a variety of individual components. Each of these components may need to be assigned an HTS code when the transmission system is imported. Assigning and tracking individual HTS codes manually is an extremely complicated and difficult task, particularly for major corporations involved in significant import or export activity. Also, classification codes assigned by a corporation, which a broker may then use to classify imported goods, may be initially incorrect in the corporation's records. This may occur, for example, because of subsequent information received from third parties relating to a particular part. Information obtained from various entities within the importation process, recent U.S. Customs Rulings, and the like may bear on whether initial classifications established by corporate analysts were correct in the first instance or have been modified subsequent to their establishment. Because of these and other complexities in the classification process, such subsequent information often gets lost or overlooked.\nMoreover, complications can be inherent to the import process itself. At least four distinct phases of operation may exist as part of this process—namely, (i) a pre-entry or classification process, (ii) an entry process, (iii) a post-entry process (including audits and payment balancing), and (iv) an amendment process. Each such phase often includes a complicated set of procedures, many of which are dependent upon or interrelated to other procedures. The phases also usually include a complex set of data concerning importation procedures, which data may contain many other relevant dependencies and interrelationships. These procedures and interrelationships must often be integrated and managed in a meaningful way so as to ensure full compliance with U.S. import regulations.\nManaging and communicating necessary information to and from an importer and a customs broker can also be important. Customs brokers may be appointed by the importer to prepare the necessary paperwork for a given shipment or set of shipments, such as a U.S. Customs 7501 form, a commercial invoice, and a shipment manifest and to present those forms to U.S. Customs when the goods are imported into the country. Customs brokers may also tender tariff payments to U.S. Customs, and are subsequently reimbursed by the importer. The reimbursement process is often slow because of its cumbersome nature. Typically, reimbursement requests or invoices are submitted by the broker to the importer, who must then enter data from the broker's request into a computerized system. Then a process of invoice balancing, reimbursement approval and generating payment for the broker begins. Each step of this process may require a different person or different data, and there is no centralized system for providing the data or access to multiple persons or entities.\nBecause of the difficulties described above, and other complexities involved in the classification, reporting and customs broker reimbursement process, broker reimbursement may be unduly cumbersome and time consuming."}
{"text": "1. Field of the Invention\nThe present invention relates to a wire harness-fixing clamp for holding a wire harness comprising wires or a cable, and more particularly to such a clamp suited for holding a self-supporting-type communication cable called xe2x80x9ca hanger gourd-shaped cablexe2x80x9d.\nThe present application is based on Japanese Patent Application No. Hei. 11-233535, which is incorporated herein by reference.\n2. Description of the Related Art\nFIG. 3 is a cross-sectional view showing a self-supporting-type cable 1 (called xe2x80x9ca hanger gourd-shaped cablexe2x80x9d) which is a kind of communication cable.\nThe self-supporting-type cable 1, having the illustrated configuration, includes a cable body portion 2, and a support wire portion 3 integrally formed with the cable body portion 2. Such self-supporting-type cables have various characteristics, and therefore have been extensively used in accordance with the purpose of use and a place where it is used. The cable body portion 2 comprises a conductor (cable core) 2a, an insulating layer 2b formed on the conductor 2a to cover the same, and a protective sheath 4 formed on the insulating layer 2b to cover the same. The support wire portion 3 comprises a support wire 3a, and the sheath 4 formed on the support wire 3a to cover the same, so that the support wire portion 3 is integrally connected to the cable body portion 2. The cable body portion 2 and the support wire portion 3 are interconnected by a sheath connection portion 4a at the boundary between the two portions 2 and 3.\nIn the installation of the self-supporting-type cable 1 of the above configuration, this cable is often fixed by clip-type clamps at suitable portions thereof spaced along the entire length thereof. One example of such clamps is a wire harness clamp disclosed in Unexamined Japanese Utility Model Publication No. Sho. 63-179717.\nThis clamp 5, shown in FIGS. 4A and 4B, includes a pair of left and right elastic holding portions 5a and 5b extending upwardly from a clamp body molded of a resin, and a holding recess 5c of an upwardly-open U-shape (in the drawings) is formed between the two elastic holding portions 5a and 5b. A retaining projection 5d is formed at an upper end of one of the elastic holding portions 5a, and bulges to be exposed to the holding recess 5c. \nWhen installing the self-supporting-type cable 1, this cable 1 is fixed to a mounting panel or the like at suitable portions thereof by the use of the clamps 5. More specifically, when the clamp 5 is mounted in an upwardly-directed posture as shown in FIG. 4A, the self-supporting-type cable 1 is usually fitted into the holding recess 5c, with the support wire portion 3 directed downwardly, so that the cable body portion 2 is disposed at a level above the support wire portion 3. The cable body portion 2, thus disposed at the higher level, is held between the left and right elastic holding portions 5a and 5b, and is retained by the elastic force of the resin, and the retaining projection 5d is press-contacted with the cable body portion 2 to prevent the same from withdrawal from the holding recess 5c. \nWhen the self-supporting-type cable 1 is set in an inverted posture in the related clip-type clamp 5 as shown in FIG. 4B, there is encountered the following disadvantage. More specifically, when the self-supporting-type cable 1 is attached to the clamp 5 in such a manner that the cable body portion 2 is directed downwardly, and is first fitted into the holding recess 5c, the support wire portion 3 projects outwardly from the clamp 5, so that the fixing of the cable is unstable. As a result, in use, the cable can flutter to produce noises, and in the worst case, the self-supporting-type cable is disengaged from the clamp 5.\nIt is therefore an object of the present invention to provide a wire harness-fixing clamp particularly best suited for holding a self-supporting-type communication cable (called xe2x80x9ca hanger gourd-shaped cable) among wire harnesses in a stable condition.\nTo achieve the above object, according to the first aspect of the present invention, there is provided a wire harness-fixing clamp which comprises a clamp body, a pair of elastic holding portions extending from the clamp body, the elastic holding portions being opposed to each other to form therebetween a holding recess, into which a wire is insertable, and which is operative to hold the wire, and a plurality of retaining projections formed respectively on the elastic holding portions to project into the holding recess, wherein when the wire is inserted into the holding recess, at least one of the retaining projections is operative to retain the wire in the holding recess. When the wire is held in the holding recess, the retaining projections are operative to engage the wire from such a direction as to prevent withdrawal of the wire from the holding recess. The retaining projections may be formed in a bulged manner respectively on inner surfaces of the elastic holding portions.\nIn the related art, the retaining projection is formed on only one of the elastic holding portions, and therefore there were occasions when the fixing of the wire harness was unstable. In the above construction of the present invention, however, the retaining projections are formed on the opposed elastic holding portions, respectively, and therefore even when the wire harness is set in any posture, the clamp can hold the wire harness in a stable manner.\nAccording to the second aspect of the present invention, the wire may include a body portion, a connection portion, and a support wire portion connected to the body portion through the connection portion. In this case, when the wire is inserted into the holding recess so that the support wire portion is first introduced into the holding recess, the body portion is engaged with the retaining projections to be prevented from withdrawal from the holding recess. Further, when the wire is inserted into the holding recess so that the body portion is first introduced and pressed into the holding recess, the body portion is engaged with one of the retaining projections and a bottom of the holding recess. The wire may be a self-supporting-type cable having a hanger gourd-shape.\nAccordingly, even when the wire is set in a proper posture or an inverted posture, with the support wire portion or the body portion first introduced into the holding recess, the body portion is positively held in a press-contacted manner, and therefore the self-supporting-type cable is positively held in a stable condition.\nAccording to the third aspect of the present invention, it is preferable that the retaining projections are offset from each other in a longitudinal direction of the holding recess. More specifically, one of the retaining projections may be formed on an intermediate portion of one of the elastic holding portions, and the other one of the retaining projections may be formed on a distal end portion of the other one of the elastic holding portions. In this case, the retaining projections press-contact with the body portion of the wire, or the one of the retaining projections and the bottom of the holding recess press-contact with the body portion of the wire.\nThat is, the retaining projections are formed on the suitable portions of the opposed elastic holding portions, respectively. More specifically, the one retaining projection is formed on one of the two elastic holding portions intermediate the opposite ends thereof, whereas the other retaining projection is formed on the distal end portion of the other elastic holding portion. With this construction, the body portion is held between the two retaining projections, and are press-contacted with these retaining projections to be prevented from withdrawal from the holding recess. Alternatively, the body portion is held between the one retaining projection and the bottom of the holding recess in a press-contacted manner, and therefore is stably held against withdrawal from the holding recess."}
{"text": "In many memory devices, including random access memory (RAM) devices, data is typically accessed by supplying an address to an array of memory cells and then reading data from the memory cells that reside at the supplied address. However, in content addressable memory (CAM) devices, data within a CAM array is not accessed by initially supplying an address, but rather by initially applying data (e.g., search words) to the array and then performing a search operation to identify one or more entries within the CAM array that contain data equivalent to the applied data and thereby represent a “match” condition. In this manner, data is accessed according to its content rather than its address. Upon completion of the search operation, the identified location(s) containing the equivalent data is typically encoded to provide an address (e.g., block address+row address within a block) at which the matching entry is located. If multiple matching entries are identified in response to the search operation, then local priority encoding operations may be performed to identify a location of a best or highest priority matching entry. Such priority encoding operations frequently utilize the relative physical locations of multiple matching entries within the CAM array to identify a highest priority matching entry. An exemplary CAM device that utilizes a priority encoder to identify a highest priority matching entry is disclosed in commonly assigned U.S. Pat. No. 6,370,613 to Diede et al., entitled “Content Addressable Memory with Longest Match Detect,” the disclosure of which is hereby incorporated herein by reference. Additional CAM devices are described in U.S. Pat. Nos. 5,706,224, 5,852,569 and 5,964,857 to Srinivasan et al. and in U.S. Pat. Nos. 6,101,116, 6,256,216 and 6,128,207 to Lien et al., the disclosures of which are hereby incorporated herein by reference.\nCAM cells are frequently configured as binary CAM cells that store only data bits (as “1” or “0” logic values) or as ternary CAM cells that store data bits and mask bits. As will be understood by those skilled in the art, when a mask bit within a ternary CAM cell is inactive (e.g., set to a logic 1 value), the ternary CAM cell may operate as a conventional binary CAM cell storing an “unmasked” data bit. When the mask bit is active (e.g., set to a logic 0 value), the ternary CAM cell is treated as storing a “don't care” (X) value, which means that all compare operations performed on the actively masked ternary CAM cell will result in a cell match condition. Thus, if a logic 0 data bit is applied to a ternary CAM cell storing an active mask bit and a logic 1 data bit, the compare operation will indicate a cell match condition. A cell match condition will also be indicated if a logic 1 data bit is applied to a ternary CAM cell storing an active mask bit and a logic 0 data bit. Accordingly, if a data word of length N, where N is an integer, is applied to a ternary CAM array having a plurality of entries therein of logical width N, then a compare operation will yield one or more match conditions whenever all the unmasked data bits of an entry in the ternary CAM array are identical to the corresponding data bits of the applied search word. This means that if the applied search word equals {1011}, the following entries will result in a match condition in a CAM comprising ternary CAM cells: {1011}, {X011}, {1X11}, {10X1}, {101X}, {XX11}, {1XX1}, . . . , {1XXX}, {XXXX}.\nApplications using CAM devices include database management, disk caching, pattern and image recognition and artificial intelligence. CAM devices are also well suited for use in routing network traffic, such as in network address lookup or packet switching. A network switch comprising a CAM device having entries therein arranged in sectors is illustrated as FIG. 1 of U.S. application Ser. No. 09/962,737, entitled “Content Addressable Memory (CAM) Devices That Can Identify Highest Priority Matches in Non-Sectored CAM Arrays and Methods of Operating Same, filed Sep. 25, 2001, the disclosure of which is hereby incorporated herein by reference. Each of the illustrated sectors is organized to contain only entries having the same number of actively masked bits, with the number of masked bits identifying entries of same priority.\nFIG. 1 herein illustrates a conventional CAM device having a plurality of CAM arrays therein arranged in a plurality of rows and columns. The CAM arrays in the first, second, third and fourth rows are illustrated as CAM00–CAM07, CAM10–CAM17, CAM20–CAM27 and CAM30–CAM37. A respective row priority encoder is also provided between each pair of CAM arrays. Thus, as illustrated, the CAM device of FIG. 1 includes sixteen (16) row priority encoders (shown as Row Priority Encoder00-Row Priority Encoder33). These row priority encoders perform final encoding of all match information generated by a respective pair of CAM arrays. A respective global word line decoder is also provided for each row of CAM arrays. As will be understood by those skilled in the art, each global word line decoder provides word line signals on global word lines to the CAM arrays of a respective row during reading and writing operations. Unfortunately, the CAM device of FIG. 1 may not achieve sufficiently high integration or sufficiently low power consumption because each of the row priority encoders typically consumes substantial chip area and requires substantial duplication of circuitry. Moreover, the pitch between global word lines that span across each row of CAM arrays may not be sufficiently large to achieve high yield and reliability when CAM arrays having many rows of normal and redundant CAM cells are utilized.\nU.S. Pat. No. 6,307,767 to Fuh also discloses a CAM device having a plurality of CAM arrays therein that are electrically coupled to a central priority encoder. In particular, FIG. 3 of the '767 patent illustrates a prior art CAM device having a central priority encoder 120 that receives match control signals from multiple CAM arrays 101–116 during a lookup operation and then, in response, generates an output address of a highest priority matching entry. FIG. 4 of the '767 patent discloses a CAM system having a plurality of CAM arrays therein that are assigned different priority levels. Circuitry is provided for identifying which of the plurality of CAM arrays has one or more matching entries of highest priority and then latching match control signals from the identified CAM array. These latched signals are provided to a respective priority encoder, which generates a respective address of a highest priority matching entry in the identified CAM array."}
{"text": "1. Field of the Invention\nThis invention relates to electronic circuits, and more particularly, to methods and mechanisms to record and/or compensate for the aging of the electronic circuits.\n2. Description of the Related Art\nOver the life an electronic circuit, the effects of aging may have an impact on its operation. Factors such as operating time, voltage, and temperature may change one or more characteristics of various circuit elements. For example, the threshold voltage of a transistor may change over the operating life of an integrated circuit (IC) in which it is implemented. A change in the threshold voltage of one or more transistors may in turn require a change in the supply voltage supplied to the IC. Generally speaking, an increase in the absolute value of a threshold voltage of one or more transistors in an IC may correspond to an increase in the required supply voltage for correct operation.\nOne common degradation mechanism that manifests itself during the aging of transistors in electronic circuits is negative bias temperature instability (NBTI), which may apply to PMOS (p-channel metal oxide semiconductor) transistors. The affects on a PMOS transistor of NBTI over a period of time may cause an increase in the absolute value of the threshold voltage, along with a decrease in drain current and transconductance. As the absolute value of the transistor's threshold voltage increase over time, a higher supply voltage value is required to ensure that the circuit operates properly. A similar phenomenon, positive bias temperature instability (PBTI) may affect NMOS (n-channel metal oxide semiconductor) transistors. Another potential degradation mechanism is hot carrier injection (HCI), wherein electrons or holes may gain sufficient kinetic energy to overcome potential barriers between different portions of the silicon (e.g., the potential barrier between the silicon substrate and the gate dielectric). Over time, HCI may degrade the gate dielectric of a transistor, increase its sub-threshold leakage current, and may also shift the threshold voltage."}
{"text": "The following is not an admission that anything discussed below is part of the prior art or part of the common general knowledge of a person skilled in the art.\nVarious types of surface cleaning apparatus are known. Surface cleaning apparatus include vacuum cleaners. Currently, a vacuum cleaner typically uses at least one cyclonic cleaning stage. More recently, cyclonic hand vacuum cleaners have been developed. See for example, U.S. Pat. No. 7,931,716 and US 2010/0229328. Each of these discloses a hand vacuum cleaner which includes a cyclonic cleaning stage. U.S. Pat. No. 7,931,716 discloses a cyclonic cleaning stage utilizing two cyclonic cleaning stages wherein both cyclonic stages have cyclone axis of rotation that extends vertically. US 2010/0229328 discloses a cyclonic hand vacuum cleaner wherein the cyclone axis of rotation extends horizontally and is co-axial with the suction motor. In addition, hand carriable cyclonic vacuum cleaners are also known (see U.S. Pat. No. 8,146,201 and U.S. Pat. No. 8,549,703)."}
{"text": "1. Field of the Invention\nThe present invention relates to a face image obtaining apparatus for obtaining a face image to be attached to a personal paper or the like belonging to an individual, in particular, ID card, magnetic card, or the like. More specifically, the present invention relates to a face image obtaining apparatus having hand-related biological information obtaining function, as well as providing a face image of the user.\n2. Description of the Related Art\nCurrently, unattended face image obtaining apparatuses for providing face images of the users are installed on the street. Such apparatuses provide recorded face images to the users by printing on plain papers or stickers.\nIn addition, a face image obtaining apparatus capable of obtaining biological information of the user such as fingerprint and the like, as well as face image, is proposed as described, for example, in International Patent Publication No. WO2005/050508. In the apparatus disclosed in the aforementioned patent publication, a face image of the user is recorded first, then the biological information. The apparatus records monitoring photographs, including a face image of the user, before and after obtaining biological information in order to provide a proof record when the user is switched for counterfeiting purpose. Such monitoring photographs may have a deterrent effect on the counterfeiting user switching. But, it is difficult to prevent such counterfeiting user switching at the site where the biological information and face image are obtained.\nIn order to prevent such counterfeiting user switching, a face image obtaining apparatus in which the face image and biological information are obtained at the same time has also been considered. If, for example, a fingerprint is obtained as the biological information, it is not an easy task for the user not accustomed to taking a fingerprint to obtain a face image and a proper fingerprint applicable to fingerprint authentication at the same time. When a user initially failed to obtain a fingerprint, if reacquisition of the fingerprint is authorized, the counterfeiting user switching may not be prevented.\nIn the mean time, another face image obtaining apparatus is also proposed as described, for example, in International Patent Publication No. WO2005/050508. In the apparatus, biological information is obtained before a face image, and when obtaining the face image, the biological information is obtained again to verify the identity of the user.\nThe face image obtaining means described in Japanese Unexamined Patent Publication No. 2005-141429 may prevent the counterfeiting user switching, but has a problem that the biological information needs to be reacquired when recording the face image, which increases the burden on the user.\nThe present invention has been developed in view of the problem described above, and it is an object of the present invention to provide a face image obtaining apparatus capable of reliably obtaining biological information applicable to authentication, and preventing counterfeiting user switching without increasing the burden on the user."}
{"text": "This invention relates to an automotive vehicle, especially a passenger car, having a bumper, at least two radiators, and an air duct for guiding the cooling air stream over the radiators. The air duct has an inlet port provided with air directing means.\nGerman Unexamined Published Patent Application 2,306,317 shows an automotive vehicle which has two series-arranged radiators extending transversely to the driving direction. A common air duct is associated with the two radiators. The air duct includes two walls forming a tunnel and air directing means provided between the walls in the area of the cooling air inlet port. The air directing means is formed from a separately manufactured synthetic resin component.\nA disadvantage of this arrangement is that the radiators, disposed in a common duct, are not exposed to an adequately defined throughflow of cooling air, as is required, for example, for a supercharger air cooler and an engine radiator.\nMoreover, this air duct must be large in cross-sectional area to accomodate the necessary air flow for the radiators, since the cooling air inlet port is located in an area of relatively low dynamic pressure. Furthermore, design freedom in the front end space is impaired by the large-area structure of the air duct. Another drawback is that the air duct, which consists of several parts, is relatively expensive to manufacture and assemble.\nIt is an object of this invention to provide, for an automotive vehicle having at least two radiators, a cooling air guide system which is simple in construction, which functions satisfactorily, and which can be integrated into the vehicle with few problems.\nThese objects are attained in a cooling air guidance system which comprises a separate air duct for each radiator. Each air duct has an inlet port provided with means for directing the cooling air flow. The inlet ports and the air directing means are arranged in an elastic cover of the bumper. One of the air ducts extends, in part, through an opening in a rigid support member for the bumper. To increase cooling efficiency, at least one of the inlet ports is located in an area of maximum dynamic pressure, e.g., below a nose-shaped projection of the bumper cover.\nThe two radiators and their respective air ducts are preferably arranged in upper and lower vertical relation. Each air duct is formed by an air guide housing, detachably connected to the radiator and the inlet port. The housings are preferably made of a synthetic resin by a blow-molding technique.\nThe lower inlet port is formed by a plurality of vertically spaced slotted openings which extend transversely to the vehicle. Ribs in the bumper cover extend between the slotted openings and serve as means for directing the cooling air stream.\nThe bumper is supported by a rigid support member which is open, relative to the longitudinal axis of the vehicle, and which is situated directly behind the upper inlet port. The support member forms a portion of the upper air duct. The support is formed by two generally horizontal plate members, spaced vertically a distance from each other. The panels are held in position by locally arranged holders and spacer elements. The panels are fastened to the holders and spacers by fasteners, such as screws, rivets and the like.\nThe primary advantage attained by the invention is improved cooling as a result of providing each radiator with its own air duct. By locating the inlet ports in an area of maximum dynamic pressure, an increased air flow rate is attained with a smaller air duct cross section. This allows for reductions in the cross section of the air duct and the dimensions of the radiators. The inlet ports can also be made substantially smaller in size if they are located in an area of high, rather than low, dynamic pressure. Since the bumper (support and cover) is utilized for forming the cooling air guide system, the space available at that location is more effectively used.\nBy arranging the inlet ports and the air directing means in the elastic cover of the bumper and by manufacturing the air guide housings of a synthetic resin by a blow-molding method, a simple and economically producible cooling air guide system is created which can be mounted in a rapid and simple fashion. The bumper support, which consists of two panels with associated holder and spacer elements, is simple in construction and can be inexpensively manufactured.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "From DE 102 60 000 B4 by the present applicant a wheel-force dynamometer for measuring tire forces is known, wherein a vehicle wheel is fixed onto a wheel axle which is mounted by roller bearings in a hollow shaft. The hollow shaft is mounted hydrostatically in a housing fixed to a frame and has a collar in which force sensors for the measurement of forces and torques are arranged. The forces acting on the wheel are thus transmitted, via the wheel axle, to the hollow shaft, which for its part “floats” in a frictionless manner by hydrostatic means in the housing. During the measurement of tire forces by a wheel-force dynamometer measurement errors can occur, which are determined by the design of the measuring device and its vibration behavior."}
{"text": "The invention relates to an inductive sensor head for detecting a ferrous, ferric and/or nonferrous electrically conducting objects buried in a surrounding medium.\nMetal detectors for detecting ferrous or non-ferrous objects in media like walls of concrete, brick, plaster or the like or in the ground based upon the disturbance or modulation of the inductive coupling between two coils are known in the art. For example, U.S. Pat. No. 5,729,143 describes a microprocessor controlled metal detector which uses a transmitter coil providing a periodically varying magnetic field in combination with a receiver coil connected thereto in an inductive bridge. The detector comprises means for automatically balancing the two overlappingly arranged coils and electronically compensating any initial coil misalignments or unwanted signals, in particular, during an initial calibration step. In a known metal detector, one of the coils, the field coil, generates an alternating magnetic field while the other coil, the sense coil, measures changes caused by a ferrous or non-ferrous material coming into the magnetic flux field while moving the detector over the medium containing the hidden disturbing object.\nA problem with the known metal detectors is, on the one hand, the relatively large size, which is unavoidable due to the side-by-side arrangement of the field coil and the sensor coil and, on the other hand, the fact that the detector must be swept over a certain search area in a kind of scanning process.\nIt is an object of the invention to provide an inductive sensor head which is small in size and may be used as a hand-held tool or may be integrated into an electric hand-held tool, preferably, a drill hammer.\nIt is a further object of the invention to provide an inductive sensor head which provides sufficient clear information about a hidden ferrous or non-ferrous electrically conducting object without the necessity of sweeping the sensor head over a certain working area of the medium in which said object may be buried.\nThe invention provides an inductive sensor head for detecting of ferrous or non-ferrous electrically conducting objects hidden. In particular, such a sensor head comprises at least one larger diameter field coil with a small axial length-to-diameter-ratio and at least one twin pair of co-axially arranged sense coils both having a small diameter compared to the diameter of the field coil. Preferably, the inductance of the sense coil is significantly higher than the inductance of the field coil. The higher the inductance the more sensitive the sense coil is to magnetic changes and the less gain is needed in the amplifiers that follow such elements. The common axis of the twin pair of sense coils extends perpendicular to the axis and in a diameter direction to the field coil, and the axis is positioned in a plane of the winding plane of the field coil or in a plane essentially parallel to the winding plane of that field coil. Further, the two sense coils are positioned in an equal distance from the center of the field coil such that they are penetrated by the same magnetic flux direction of the flux field emanating from the field coil when excitated by an electric current.\nFor achieving better positional information, in particular for resolving depth information in relation to a hidden object, e.g., a reenforcing bar (xe2x80x9crebarxe2x80x9d in the following) from a single position measurement cycle, a significant improvement of the invention is achieved if a twin pair of coaxially positioned field coils is provided. The mutual axial distance of the two field coils can be rather close and may preferably be less than their internal diameter. As a rule, the distance between the field coils is arranged such that the difference in magnetic field strength on a rebar is sufficiently large that it can be accurately measured. In addition, two twin pairs of sense coils with orthogonally arranged axes are positioned in a center plane parallel and approximately at a halfway distance between the winding planes of the two field coils.\nAs will be described in the following further details, the invention also provides an advantageous driving circuit for the combination of a twin pair of field coils and a double twin-set of sense coils, wherein additional correction coils are provided in series connection with each of the two field coils in order to minimize magnetic offsets due to the fact that the sense coils cannot be or are difficult to be exactly positioned in the magnetic null position of both field coils.\nThe various features of novelty which characterize the invention are pointed out with particularity in the claims annexed to and forming a part of this disclosure. For a better understanding of the invention, its operating advantages and specific objects attained by it use, references should be had to the drawings and description matter in which there are illustrated and described preferred embodiments of the invention."}
{"text": "The present invention relates to a rotor for a rotary electric machine.\nA structure of a conventional rotor core of a wound rotor having a rotor coil for a rotary electric machine is constructed by punching a silicon steel plate in a disc shape and laminating the disc-shaped plates.\nHowever, since the punching die used to punch the steel plate is large in size in the case of a large-sized rotary electric machine, a rotor core is constructed, due to the above mentioned restriction in the manufacture of the rotor core, for example, as shown in FIG. 4, by punching a silicon steel plate in a sector shape, then arranging sector-shaped steel plates in a circle as shown in FIGS. 6(A) and 6(B), and laminating them.\nIn FIGS. 4, 5, 6(A) and 6(B), reference numeral 1 denotes a rotor core, symbol la denotes rotor core pieces made of a punched silicon steel plate, numeral 2 denotes slots to which a rotor coil is inserted, numeral 3 denotes a clamping bolt hole through which a clamping bolt is penetrated, numeral 4 denotes keyways for coupling the rotor core 1, for example, to a spider boss, numeral 5 denotes a ventilation duct for ventilating to cool the rotor, numeral 6 denotes a rotor coil, numeral 7 denotes a retainer for clamping the rotor core to integrate them by clamping the rotor core pieces la, numeral 8 denotes a ventilation fan for blowing cooling air to the stator, numeral 9 denotes a spider boss, and numeral 10 denotes a rotational shaft. The slots 2 are shown only in FIG. 4 for the simplification of the illustration.\nThen, a method of manufacturing the conventional rotor will be described.\nThe rotor core 1 having a large diameter is generally formed by punching a thin silicon steel plate in a sector shape as shown in FIG. 4 to form rotor core pieces 1a. Then, the rotor core pieces 1a are stacked in a cylindrical shape while providing the ventilation duct 5 as shown in the sectional view of FIG. 5. In other words, a plurality of the rotor core pieces 1a are disposed horizontally to form one circular shape, and then the layers each formed of the one circular-shaped thin plates made of the rotor core pieces 1a are sequentially stacked upward similarly to the above. In the case of FIG. 4, 10 sheets of the core pieces 1a are disposed horizontally in one circular shape as one layer, and the layers thus formed are sequentially stacked upward. As described above, the rotor core pieces 1a thus stacked are clamped by penetrating the clamping bolts through the clamping bolt holes 3, and the rotor core 1 is constructed by providing the retainers 7 for clamping the rotor core at the upper and lower ends of the rotor core pieces 1a thus stacked, penetrating the clamping bolts therethrough, and clamping them with the clamping bolts in the axial direction.\nThe slots 2 for inserting the rotor coil 6 are formed at the outer periphery of the rotor core 1. After the rotor core 1 is constructed, the rotor coil 6 is inserted into the slots 2, and fixed in the slots 2 by inserting wedges to the outer peripheral ends of the slots 2.\nA torque is transmitted between the spider boss 9 and the rotor core 1 through keys inserted into the keyways 4.\nAs described above, the rotor core pieces 1a are arranged horizontally in one circular shape, and the layers of the circular-shaped thin plates made of the rotor core pieces 1a are sequentially stacked. FIGS. 6(A) and 6(B) show the case where the rotor core pieces 1a are lap stacked at each 1/4 of one circle. More specifically, the rotor core pieces 1a shown in FIG. 4 are aligned in one circular shape on the outer periphery of the spider boss 9. In the case of the rotor core piece shown in FIG. 4, when ten sheets of the rotor core pieces 1a are aligned horizontally, they form one circular shape, and the split positions of the sectors are as designated by thick solid lines indicated by numeral 11 in FIG. 6(A).\nWhen the second layer of the rotor core pieces 1a shown in FIG. 6(B) is so arranged in one circular shape that the split position is disposed at the position displaced by 1/4 of one sector from the split position 1 of the first layer of the rotor core pieces 1a, the split position of the second sector is disposed at the portion designated by a fine solid line indicated by numeral 12 in FIG. 6(A).\nSimilarly, when the third and fourth layers of the rotor core pieces 1a are stacked by displacing them by 1/4 of one sector from the split position, the split positions of the sectors are disposed at the positions designated by a broken line indicated by numeral 13 and a dotted broken line indicated by numeral 14, respectively.\nThe conventional rotor core 1 is constructed in the cylindrical structure by sequentially stacking the layers of the rotor core pieces 1a arranged horizontally in one circular shape as described above, providing the ventilation ducts 5 between several layers, and stacking the layers of the rotor core pieces 1a to a predetermined core length.\nThe rotor core 1 is constructed by stacking the layers of the rotor core pieces 1a, and then integrally clamping the stacked rotor core pieces through the retainers 7 disposed at the upper and lower ends thereof with the clamping bolts.\nSince the conventional rotor core is constructed as described above, it is necessary to disassemble the rotor coil and to disassemble all the rotor core pieces after the rotor is assembled and tested for its factory test if the entire rotor cannot be transported as a unit due to limitations in transportation. Thus, it is also needed to reassemble the rotor coil and the rotor core pieces at a site after the disassembled rotor is transported to a place where the rotor is to be installed. Therefore, it takes a long time to disassemble the rotor in the factory, to transport the disassembled rotor and to reassemble the rotor, and there may be a problem in reliability of the performance of the reassembled rotor. It is an object of this invention to overcome this problem."}
{"text": "Retail, office and home environments frequently make use of shelving, and lighting to illuminate or even showcase items on a shelf. Movable shelves are often supported on shelf brackets, which attached to shelf support columns. Typically, each shelf support column attaches to a wall or other supporting member, and has a series of apertures to which the shelf brackets affix. A shelf or a series of shelves is generally supported by two shelf support columns, one at each end of the shelf, and optionally one or more additional shelf support columns in the middle. It would be advantageous to provide for lighting or other electrical needs at the shelves.\nU.S. Pat. No. 5,695,261 describes a modular furniture system having vertical standards with internal conductors. The conductors are coupled to a source of low-voltage electrical power such that adjacent vertical standards are of differing electrical polarity. Brackets supporting the shelves of the modular furniture system make electrical contact with the internal conductors when a bracket is engaged in a vertical standard. A light fixture, coupled between a pair of brackets engaged with adjacent vertical standards, can be energized. Yet, there is a need in the art for improvements and alternatives to the above-described system."}
{"text": "The present invention relates to a composition comprising a polyarylene ether and a dispersible reactive solvent. Polyarylene ethers (PAEs) are a class of thermoplastic resins with excellent mechanical and electrical properties, heat resistance, flame retardancy, low moisture absorption, and dimensional stability. These resins are widely used in automobile interiors, particularly instrument panels, and electrical as well as electronic applications.\nPAEs are very difficult to process (for example, by injection molding) as a result of their high melt viscosities and their high processing temperature relative to their oxidative degradation temperature. Consequently, PAEs are commonly blended with compatible polymers such as polystyrene (WO 97/21771 and U.S. Pat. No. 4,804,712); polyamides (U.S. Pat. No. 3,379,792); polyolefins (U.S. Pat. No. 3,351,851); rubber-modified styrene resins (U.S. Pat. Nos. 3,383,435 and 3,959,211, and Ger. Offen. No. 2,047,613); and mixtures of polystyrene and polycarbonate (U.S. Pat. Nos. 3,933,941 and 4,446,278). Unfortunately, improvements in processing have generally been obtained at the expense of flexural modulus, flexural strength, or heat distortion temperature.\nEpoxy resins have also been investigated as a reactive solvent for the PAE. (See Venderbosch, R. W., \"Processing of Intractable Polymers using Reactive Solvents,\" Ph.D. Thesis, Eindhoven (1995); Vanderbosch et al., Polymer, Vol. 35, p. 4349 (1994); Vanderbosch et al., Polymer, Vol. 36, p. 1167 (1995a); and Vanderbosch et al., Polymer, Vol. 36, p. 2903 (1995b)). In this instance, the PAE is first dissolved in an epoxy resin to form a solution that is preferably homogeneous. An article is then shaped from the solution, and the solution is cured at elevated temperatures, resulting in a phase separation that can give a continuous PAE phase with epoxy domains interspersed therein. The properties of the finished article are primarily determined by the PAE; however, the use of an epoxy resin as a reactive solvent for the PAE is not practical in a continuous melt process like injection molding because the epoxy resin needs a curing agent to set. The curing agent will, over time, accumulate in the injection molding barrel, thereby fouling the machine. Furthermore, the cure and subsequent phase separation has to take place at at least 150.degree. C., which is impractical in a molder environment.\nIn view of the deficiencies in the art, it would be desirable to find a reactive solvent that would solve the processing problems inherent in some reactive solvents for PAE, without deleteriously affecting the physical properties of the PAE."}
{"text": "Mobile or wireless communications networks are capable of carrying traffic (e.g., voice traffic, data traffic, etc.) between mobile stations and other endpoints, which can be another mobile station or an endpoint connected to a network such as a public-switched telephone network (PSTN) or a packet data network (e.g., local area networks, the Internet, and so forth). Examples of wireless communications networks include those that operate according to the GSM (Global System for Mobile) or UMTS (Universal Mobile Telecommunications System) standards, as defined by the Third Generation Partnership Project (3GPP). Another type of wireless communications network is according to the Code Division Multiple Access (CDMA) 2000 standards, as defined by 3GPP2.\nA media gateway is often provided between a radio access network (which is part of the core network of the wireless communications network) and a fixed network such as a PSTN, an Internet Protocol network, or other network. The media gateway performs various signal processing tasks with respect to bearer traffic (e.g., voice traffic) communicated between the radio access network and the fixed network. The signal processing includes encoding and decoding of the bearer traffic, with the encoding and decoding typically performed by a digital signal processor (DSP). The DSP is a shared resource that can be shared by multiple call sessions (or call contexts). Multiple events (e.g., encoding or decoding events) associated with the call contexts are scheduled so that the DSP processes the events one at a time according to the schedule. Conventional scheduling techniques do not allow for efficient handling of dynamically changing numbers of events that are to be processed by the DSP."}
{"text": "With portable electric tools there is a need to simplify assembly to both reduce production costs and to reduce the risk of assembly errors. This has become more important as such tools have become more sophisticated in their functioning.\nIn the manufacture of electric motors for such tools, it is becoming increasingly common practice to wind the field coils mechanically on to the stator and to provide terminations on the latter for receiving the ends of the field coil windings and which facilitate electrical connection of the windings to the commutator brushes. The stator assembly can be formed by a stack of field laminations and a plurality of coils, and be adapted for automatic connection of the coils to terminal means mounted on the stack wherein the terminal means and mounting means lie entirely within an area defined by the outline of the field laminations. Such an arrangement is disclosed in U.S. Pat. No. 4,071,793 which is hereby incorporated by reference.\nImprovements have been made in the manner of connecting the electric leads to the stator assembly. In one such arrangement a pair of blocks made from suitable insulating material such as a polysulphone are located in slots in the stator laminated stack, these blocks being provided with a pair of apertures for receiving a conductive terminal. Each terminal comprises a sleeve portion for engaging in the aperture and a channel portion connected to the sleeve portion by a short connecting neck. A wire to be attached is crimped in the channel portion. Such an arrangement is disclosed in British Pat. No. 1,402,591 which is hereby incorporated by reference. When this method of connecting electrical leads is used with the stator assembly referred to above, the stator assembly can be readily manufactured as a separate unit which is then easily insertable into the housing of the portable electric tool and then the electrical connections to be made to it can be made simply and effectively.\nIt has been proposed to mount a printed circuit board on a plate having attached thereto carbon brush assemblies, with the plate being attached to the housing of the tool. The armature of the electric motor passes through central openings in both the plate and the printed circuit board.\nIn order to reverse the rotational direction of drive of an electric tool, a separate reversing switch can be incorporated. However, with many forms of motors, for example, universal motors, damage can occur if the reversing switch is operated to reverse the direction of electrical supply to the motor whilst it is still rotating. To eliminate this danger of damage occurring to the electrical motor, it has been proposed to incorporate the reversing switch in a trigger switch for energizing the tool. The trigger switch is mounted, as well known, in the handle of the tool, and the actuating member of the reversing switch is disposed immediately above the trigger of the trigger switch and just below the motor compartment of the tool. The actuating member of the reversing switch and the trigger are mechanically related so that the trigger remains inoperative, i.e. it cannot be moved, until the actuating member of the reversing switch is positioned to one side of the trigger to allow the motor to be energized to rotate in one direction, or until the actuating member is positioned to the other side of the trigger to reverse the direction of rotation of the motor.\nA disadvantage of this reversing switch and trigger switch combination is that it complicates the number of electrical wires that have to feed from the handle of the tool through to the motor compartment and also the number of electrical connections that have to be made to the combined switches in the handle.\nThe present invention is concerned with further simplifying the assembly of portable electric tools.\nIt is an object of this invention to provide a portable electric tool having a reversing switch interrelated with a main energizing switch and being arranged so that the number of electrical wires feeding from the handle to the motor compartment can be reduced by at least two.\nIt is another object of this invention to provide a portable electric tool having a printed circuit board assembly in the motor compartment with the reversing switch being part of that assembly.\nIt is yet a further object of this invention to provide a portable electric tool having a comprehensive printed circuit board assembly in the motor compartment and being readily mounted on a stator lamination stack of the electric motor."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nFIG. 1 is a view illustrating an example of the semiconductor light emitting device proposed in U.S. Pat. No. 7,262,436. The semiconductor light emitting device includes a substrate 100, an n-type semiconductor layer 300 grown on the substrate 100, an active layer 400 grown on the n-type semiconductor layer 300, a p-type semiconductor layer 500 grown on the active layer 400, electrodes 901, 902 and 903 formed on the p-type semiconductor layer 500, while serving as reflective films, and an n-side bonding pad 800 formed on the n-type semiconductor layer 300 which has been etched and exposed. The n-type semiconductor layer 300 and the p-type semiconductor layer 500 can be of opposite conductive types. Preferably, a buffer layer (not shown) is provided between the substrate 100 and the n-type semiconductor layer 300. A chip having this structure, i.e. where all the electrodes 901, 902 and 903 and the n-side bonding pad 800 are formed on the opposite side of the substrate 100, with the electrodes 901, 902 and 903 serving as reflective films, is called a flip-chip. The electrodes 901, 902 and 903 are made up of an electrode 901 (e.g., Ag) with a high reflectance, an electrode 903 (e.g., Au) for bonding, and an electrode 902 (e.g., Ni) for preventing diffusion between materials of the electrode 901 and materials of the electrode 903. While this metal reflective film structure has a high reflectance and is advantageous for current spreading, it has a drawback that the metal absorbs light.\nFIG. 2 is a view illustrating an example of the semiconductor light emitting device proposed in JP Pub. No. 2006-120913. The semiconductor light emitting device includes a substrate 100, a buffer layer grown on the substrate 100, an n-type semiconductor layer 300 grown on the buffer layer 200, an active layer 400 grown on the n-type semiconductor layer 300, a p-type semiconductor layer 500 grown on the active layer 400, a light-transmitting conductive film 600 with a current spreading function formed on the p-type semiconductor layer 500, a p-side bonding pad 700 formed on the light-transmitting conductive film 600, and an n-side bonding pad 800 formed on the n-type semiconductor layer 300 which has been etched and exposed. Further, a DBR (Distributed Bragg Reflector) 900 and a metal reflective film 904 are provided on the light-transmitting conductive film 600. While this structure reduces light absorption by the metal reflective film 904, it has a drawback that current spreading is relatively poor, compared with the use of the electrodes 901, 902 and 903.\nFIG. 12 is a view illustrating an example of the semiconductor light emitting device proposed in JP Pub. No. 2009-164423. In the semiconductor light emitting device, a DBR 900 and a metal reflective film 904 are provided on a plurality of semiconductor layers 300, 400 and 500, a phosphor 1000 is provided on opposite side thereof. The metal reflective film 904 and an n-side bonding pad 800 are electrically connected with external electrodes 1100 and 1200. The external electrodes 1100 and 1200 can be lead frames for a package, or electrical patterns provided on the COB (Chip on Board) or PCB (Printed Circuit Board). The phosphor 1000 can be coated conformally, or can be mixed with an epoxy resin and then used to cover the external electrodes 1100 and 1200. The phosphor 1000 absorbs light that is generated in the active layer, and converts this light to a light of longer or shorter wavelength."}
{"text": "1. Field of the Invention\nThe present invention relates generally to voltage regulators, and more particularly to a frequency sensing voltage regulator that uses the system operating frequency to limit the amount of current delivered to a load, thereby regulating the variance of the supply voltage to the load.\n2. Description of the Related Art\nVoltage regulator circuits are known in which a voltage supply to a load is regulated by controlling the current supplied to the load. Typical of such prior art structures is the use of a negative feedback circuit for sensing the output voltage and/or output current which is used for comparison with a reference voltage/reference current. The difference between the output and the reference signal is used to adjust the current supplied to a load.\nThere are problems, however, with such voltage regulators. A considerable amount of power is drawn, and thus heat dissipated, because of the use of the negative feedback circuit. In addition, the negative feedback circuit decreases the response time to sharp current fluctuations. Furthermore, the comparator circuits and reference level generating circuits take up considerable layout area when the voltage regulator is incorporated in an integrated circuit (IC) structure.\nAdditional problems also occur when a voltage regulator is used to regulate the supply voltage to a synchronous device, such as a synchronous memory device, for example an SRAM. In an SRAM, an external supply voltage, Vcc, must be maintained within a predetermined level. The external supply voltage Vcc must be regulated to produce a regulated Vcc value during periods of considerable current fluctuation. For example, an SRAM load current may quickly fluctuate between microamps and milliamps during use. Such changes in the load current can cause significant variation on the regulated Vcc value, which can result in improper operation of the SRAM or possibly even damage to the SRAM.\nThus, there exists a need for a voltage regulator that is easy to implement, does not occupy significant layout area when the voltage regulator is incorporated in an integrated circuit (IC), and provides a minimal variance of the supply voltage Vcc over a wide current range."}
{"text": "In recent years, a motor wherein current paths are electronically altered by plural transistors has been widely used as a drive motor in office automation equipment and audio-visual equipment. A disk drive apparatus, such as an optical disc drive apparatus (DVD, CD, and the like) and a magnetic disk drive apparatus (HDD, FDD, and the like), includes such a motor. As an example of such a conventional motor, a motor wherein current paths are altered by PNP-type and NPN-type bipolar power transistors is disclosed on lines 16 to 31 in the first column and FIG. 34 in the specification of the U.S. Pat. No. 5,982,118.\nFIG. 35 shows a prior art motor, and the operation of the prior art motor is described below. A rotor 2011 has a field part formed by a permanent magnet. In a position detector 2041, three position sensors detect the magnetic field of the field part of the rotor 2011. In other words, the position detector 2041 produces two sets of three-phase voltage signals, that is, Kp1, Kp2 and Kp3, and Kp4, Kp5 and Kp6, in response to the three-phase output signals of the three position sensors in response to the rotation of the rotor 2011. A first distributor 2042 produces three-phase lower-side signals Mp1 Mp2 and Mp3 in response to the voltage signals Kp1, Kp2 and Kp3, thereby controlling the activation of lower-side NPN-type bipolar power transistors 2021, 2022 and 2023. A second distributor 2043 produces three-phase higher-side signals Mp4, Mp5 and Mp6 in response to the voltage signals Kp4, Kp5 and Kp6, thereby controlling the activation of the upper-side PNP-type bipolar power transistors 2025, 2026 and 2027. Accordingly, three-phase drive voltages are supplied to windings 2012, 2013 and 2014.\nIn the prior art configuration shown in FIG. 35, the position detector 2041 comprises three position sensors for detecting the rotational position of the rotor 2011. This has caused the necessity of a substantial space for installing these position sensors and the complexity of the wiring, resulting in an increase in cost. On the other hand, a motor with no position sensor is disclosed on line 54 of the second column to line 45 of the third column and FIG. 1 in the specification of the U.S. Pat. No. 5,473,232. In the motor, the back-electromotive forces of the windings are detected so as to obtain the rotational position of the rotor. In the motor with no position sensor, however, the rotational position cannot be detected correctly at a low rotational speed of the motor, since the magnitudes of the back-electromotive forces become too small to detect at a low rotational speed of the motor. So, it is difficult to drive and control the motor at a low speed. In particular, in the case when the rotational speed is controlled by using a pulse signal which responds with the detected back-electromotive forces, a large fluctuation occurs in the rotational speed of the motor at a low speed because of the inaccurate detection of the pulse signal.\nA motor with a single position sensor is disclosed on line 30 of the fifth column to line 41 of the 12th column and FIG. 1 in the specification of the U.S. Pat. No. 5,729,102. In the motor, the rotational electrical angle is estimated from the output of the single position sensor, and sinusoidal currents are supplied to the windings in response to the estimated rotational electrical angle. However, in the case of the motor with the single position sensor, positional information in the stop state of the rotor is insufficient, thereby starting and acceleration of the motor with the single position sensor becomes unstable. Accordingly, the starting and acceleration of the rotor are not carried out smoothly, resulting in a starting failure. Furthermore, in the configuration of the motor according to the specification of the U.S. Pat. No. 5,729,102, it is difficult to estimate the rotational electrical angle with a fine step resolution. In particular, the error in the estimated electrical angle becomes larger at a higher rotational speed. Accordingly, precise rotation control of the motor has been difficult.\nIn an optical disc drive apparatus for reproducing signals from DVD-ROM, CD-ROM, and CD disks, stable operation is required over a wide range of rotational speed from 10,000 rpm for high-speed reproduction to 200 rpm for CD reproduction. In a rewritable disk drive apparatus for recording information on a high-density disk and/or reproducing information from a high-density disk such as DVD-RAM/RW and CD-R/RW, the disk is required to be rotated precisely. In these disc drive apparatuses, it is necessary to smoothly start and accelerate the disk and to carry out information reproduction in a short time. In addition to the optical disc drive apparatuses, magnetic disc drive apparatuses such as HDD and FDD are also required to be low cost and to carry out stable rotation of the disk during the whole operation which includes the operation of the starting and acceleration thereof."}
{"text": "Pectin polymers are important constituents of plant cell walls. Pectin is a hetero-polysaccharide with a backbone composed of alternating homogalacturonan (smooth regions) and rhamnogalacturonan (hairy regions). The smooth regions are linear polymers of 1,4-linked alpha-D-galacturonic acid. The galacturonic acid residues can be methyl-esterified on the carboxyl group to a varying degree, usually in a non-random fashion with blocks of polygalacturonic acid being completely methyl-esterified.\nPectinases can be classified according to their preferential substrate, highly methyl-esterified pectin or low methyl-esterified pectin and polygalacturonic acid (pectate), and their reaction mechanism, beta-elimination or hydrolysis. Pectinases can be mainly endo-acting, cutting the polymer at random sites within the chain to give a mixture of oligomers, or they may be exo-acting, attacking from one end of the polymer and producing monomers or dimers. Several pectinase activities acting on the smooth regions of pectin are included in the classification of enzymes provided by the Enzyme Nomenclature (1992) such as pectate lyase (EC 4.2.2.2), pectin lyase (EC 4.2.2.10), polygalacturonase (EC 3.2.1.15), exo-polygalacturonase (EC 3.2.1.67), exo-polygalacturonate lyase (EC 4.2.2.9) and exo-poly-alpha-galacturonosidase (EC 3.2.1.82).\nPectate lyases have been cloned from different bacterial genera such as Erwinia, Pseudomonas, Klebsiella and Xanthomonas. Also from Bacillus subtilis (Nasser et al. (1993) FEBS 335:319-326) and Bacillus sp. YA-14 (Kim et al. (1994) Biosci. Biotech. Biochem. 58:947-949) cloning of a pectate lyase has been described. Purification of pectate lyases with maximum activity in the pH range of 8-10 produced by Bacillus pumilus (Dave and Vaughn (1971) J. Bacteriol. 108:166-174), B. polymyxa (Nagel and Vaughn (1961) Arch. Biochem. Biophys. 93:344-352), B. stearothermophilus (Karbassi and Vaughn (1980) Can. J. Microbiol. 26:377-384), Bacillus sp. (Hasegawa and Nagel (1966) J. Food Sci. 31:838-845) and Bacillus sp. RK9 (Kelly and Fogarty (1978) Can. J. Microbiol. 24:1164-1172) has been reported, however, no publication was found on cloning of pectate lyase encoding genes from these organisms. All the pectate lyases described require divalent cations for maximum activity, calcium ions being the most stimulatory.\nGenerally, pectinase producing microorganisms exhibit a broad range of pectin degrading or modifying enzymes. Often the microorganisms also produce cellulases and/or hemicellulases and complex multi-component enzyme preparations from such microorganisms may be difficult to optimise for various applications, they even may contain enzymes with detrimental effect. Thus, it is an object of the present invention to provide a pectin degrading enzyme exhibiting only the desired effects e.g. in detergents or different industrial processes."}
{"text": "The present invention relates generally to fluid power systems, and more particularly to a fluid power system which includes two fluid power circuits and a means for controlling the pressure in one of the circuits in response to the pressure in the other circuit. Still more particularly, the invention relates to such a fluid power system for use in a dewatering filter press.\nDewatering filter presses are used to remove excess water from a slurry. In such filter presses, water pressure acts against a diaphragm to squeeze the water from the slurry while hydraulic pressure acts against a piston to provide a clamping force holding the press closed. In this type of press, the magnitude of the clamping force, and hence the magnitude of the hydraulic pressure in the clamping piston, must be proportional to the water pressure acting against the diaphragm. In the prior art, this is accomplished by using several pressure activated switches that are responsive to the water pressure acting against the diaphragm. This arrangement produces a clamping pressure that is not a linear function of the water pressure but instead is a stair step function of the water pressure."}
{"text": "1. Field of the Invention\nThe present invention relates to a photography device.\n2. Description of the Related Art\nAmong photography devices in which a lens for photography is protruded through a protrusion aperture of a casing body, there are photography devices which are structured such that the protrusion aperture is covered by a lens barrier when the lens for photography is retracted into the casing body (stored).\nFor example, Japanese Patent No. 2,712,160 discloses a camera which is equipped with an urging mechanism and an operation mechanism. The urging mechanism urges the lens barrier in a closing direction. The operation mechanism moves the lens barrier to an open position, against urging force from the urging mechanism, and retains the lens barrier. Accordingly, operations to open and close the lens barrier are implemented simply and safely.\nHowever, in an operation to close the lens barrier, the lens barrier temporarily abuts against a peripheral surface of a lens barrel and movement toward the closed position is obstructed. Thereafter, when the barrel has moved to a camera body side end of the movement thereof, the lens barrier is closed by the urging force from the urging mechanism. In this operation, because the lens barrier abuts against the lens barrel, there is a risk that damage or breakage may occur.\nSimilarly, in a retraction-type camera disclosed in Japanese Patent No. 3,078,408, when a barrier is being moved from a completely open position to a completely closed position, a trigger switch operates and a lens barrel starts to retract. When the barrier is moved further in the closing direction, clicking means operates, prompting a user to wait until the lens barrel has completely retracted. Then, when the retraction has completely finished and it is possible for the barrier to completely close, the barrier is moved so as to pass over the clicking means, and is completely closed.\nIn this structure, although impacts between the barrier and the lens barrel can be avoided at least, movement of the barrier is temporarily stopped at an intermediate position during complete closing of the barrier. Thus, there is a problem in that operability is poor.\nJapanese Patent Application Laid-Open (JP-A) No. 7-49515 discloses a camera (lens device) which is equipped with a motor for implementing protrusion and withdrawal of a moveable barrel, and which is structured so as to rotate a lens barrier for closing while the moveable barrel is being withdrawn.\nWith this structure, a mechanism for transmitting driving force of the motor to the lens barrier is necessary, and there is a disadvantage in that the structure is complicated."}
{"text": "The present invention relates to a flexible polyurethane foam, its production method and a material system for its production.\nIn recent years, along with progress in an automobile industry, not only improvement in performance of an automobile but also a high degree of vehicle interior and improvement in interior comfort have been desired, and development of a flexible polyurethane foam for seat (hereinafter referred to as a flexible foam) having more excellent cushioning properties has been strongly desired in view of improvement in comfort to sit on and comfort to ride in.\nHeretofore, as a seat cushion, a combination of metal springs and a pad material comprising a flexible foam has been used widely. However, in recent years, a seat for an automobile so-called full-foam type, wherein a flexible foam itself is made to have spring properties so that metal springs are not used, tends to be employed, with a demand for cost saving, weight saving, etc. The full-foam type seat tends to be thick since no metal spring is used together, and properties of the flexible foam are factors which greatly contribute to comfort to sit on and comfort to ride in of the seat. Namely, importance has been given to static characteristics and dynamic characteristics as indices of comfort to sit on and comfort to ride in, in development of a flexible foam. Particularly among the static characteristics, control of the feeling of support at the initial stage of sitting and the feeling of bottom out at the terminal stage of sitting are important.\nWhen a person actually sits on a seat equipped with a pad material made of a flexible foam, the flexible foam is compressed and indented, and the position of e.g. the bottom falls to a certain height. As a method of measuring the static characteristics (static seat feeling), a load test in accordance with a performance test method of a pad material for a seat for an automobile in JASO (Japanese Automobile Standards Organization) B408-89, wherein a deflection is measured to obtain a load-deflection loop, is used. In FIG. 2 is shown a load-deflection loop of a conventional flexible foam. In the load-deflection loop, the value of the tangent of the tangent line at the pressing side (upper side in FIG. 2) at each load is the static spring constant. Namely, when the static spring constant is high, the curve has a steep slope and when the static spring constant is low, the curve has a gentle slope.\nThe load-deflection loop of a flexible foam usually rises quickly at the initial stage of application of the load, then gently increases at the inflection point 1 and with further application of the load, it rapidly increases at the inflection point 2. A flexible foam excellent in the feeling of support is considered to be such a flexible foam that the curve steeply rises at the initial stage of application of the load, and the static spring constant is high from the starting point to the vicinity of the inflection point 1. Further, a flexible foam which provides no feeling of bottom out and provides favorable comfort to sit on, is considered to be such a flexible foam that the static spring constant is low even in the high load region after the inflection point 2. Adjustment of hardness of the flexible foam has conventionally been carried out to obtain the feeling of support, however, in such a case, the static spring constant in the high load region tends to be high, and the feeling of bottom out can not be overcome.\nGenerally, a full-foam type flexible foam for seat provides a small deflection in the high load region as the static characteristics, it provides the feeling of bottom out, and the comfort to sit on tends to deteriorate, as compared with the foam for seat in combination with metal springs. In order to overcome such problems, a method has been known to make the foam thick or to increase the density to make the deflection large, however, problems of increase in cost and increase in the seat weight can not be overcome.\nJP-A-11-322875 proposes a method of adding a specific fluorine type surfactant to a material for production of a polyurethane foam composed mainly of a polyol and an isocyanate component so as to increase the deflection. However, there are such problems that the feeling of support is insufficient, the fluorine type surfactant itself is expensive, thus increasing the cost, and the surface tension of the surfactant is low, whereby the inside of a mold at the time of foam production may be polluted.\nFurther, JP-A-5-320304 proposes a method for producing a polyurethane foam, which comprises reacting a polyol component comprising a polyol, a catalyst, a blowing agent and other additives, with a polyisocyanate component, wherein a specific bifunctional secondary amine is added, to produce a foam which provides a small slope of the pressing side curve at 75% deflection as a characteristic to evaluate the feeling of bottom out, which provides no feeling of bottom out, and which has a well balance among softness, degree of falling and vibrating property. However, the feeling of support is insufficient, and durability particularly heat and humid permanent compression set is insufficient.\nNamely, such a flexible foam is not present that the static spring constant is high in the low load region, the feeling of support and rigidity are satisfied, and further, the static spring constant is low in the high load region, the deflection is large and the feeling of bottom out is overcome.\nUnder these circumstances, it is an object of the present invention to provide a flexible foam for an automobile seat, with which the feeling of bottom out is overcome, the feeling of support is favorable, and the comfort to sit on is significantly improved, and which has favorable vibrating property and durability, and a method for producing said foam and a material system for production of said foam.\nThe present inventors have conducted extensive studies to overcome the above problems and as a result, they have found a foam which exhibits a load-deflection loop which satisfies characteristic conditions in a load test of the flexible foam or in a load test of its core portion, and found that a foam having such a characteristic load-deflection loop significantly improves the comfort to sit on as the static seat feeling, when an automobile seat is molded (see FIG. 1). A foam which satisfies such conditions has such characteristics that the deflection is small at the initial stage of application of the load, the feeling of support is sufficient, and after that region, the rate of change of the static spring constant is small, and the curve gently increases. Namely, it is estimated that at the initial stage of sitting, the feeling of support is sufficient, whereby the foam is excellent in body weight support, and after a person sits on the foam, it is well bent, whereby the body weight is softly supported, and accordingly the feeling of bottom out is overcome, and the comfort to sit on becomes favorable.\nNamely, the present invention provides a flexible polyurethane foam, wherein in a load (N)-deflection (mm) loop obtained by a load test of the flexible polyurethane foam, the X value as calculated from the formula (1) is at most 4.2:\nX=(static spring constant at a load of 883N)/(static spring constant at a load of 98N)xe2x80x83xe2x80x83(1)\nThe present invention further provides a flexible polyurethane foam, wherein in a load (N)-deflection (mm) loop obtained by a load test of a core portion of the flexible polyurethane foam, the Y value as calculated from the formula (2) is at most 1.8:\nY=(static spring constant at a load of 127.4N)/(static spring constant at a load of 19.6Nxe2x80x83xe2x80x83(2)\nThe present invention further provides a method for producing the above flexible polyurethane foam, which comprises reacting a polyol having from 5 to 85 mass % of an oxyethylene/oxypropylene random chain, the oxyethylene group content in the random chain being from 3 to 40 mass % and the content of the terminal oxyethylene block chain being from 3 to 40 mass %, and having a hydroxyl value of from 10 to 56 mgKOH/g and a degree of unsaturation of at most 0.04 meq/g, with a polyisocyanate in the presence of a blowing agent and a catalyst.\nStill further, the present invention provides a material system for production of the above flexible polyurethane foam, which comprises a polyol system liquid containing a polyol having from 5 to 85 mass % of an oxyethylene/oxypropylene random chain, the oxyethylene group content in the random chain being from 3 to 40 mass % and the content of the terminal oxyethylene block chain being from 3 to 40 mass %, and having a hydroxyl value of from 10 to 56 mgKOH/g and a degree of unsaturation of at most 0.04 meq/g, and a polyisocyanate."}
{"text": "This invention relates to controlling the operation of refrigeration systems which contain multiple compressors. More particularly, the present invention is concerned with refrigeration systems of the type having multiple compressors fed from a common suction manifold and which deliver refrigerant gas under pressure to a common head pressure manifold.\nIt is known from U.S. Pat. No. 3,599,006 to John L. Harris granted Aug. 10, 1971 and entitled \"Condition Control Device and System\" to provide multiple compressors in a refrigeration system. The system is arranged so that the compressors are started in sequence, allowing each compressor to come up to speed before the next compressor is started. The system also interposes a delay between the stopping and restarting of the respective compressors so that a compressor is not started under heavy load. The system is essentially an electromechanical system responsive to analog signals and completely analog in nature.\nIt is known from U.S. Pat. No. 4,128,854 to Robert T. Ruminsky issued Dec. 5, 1978 and entitled \"Compressor Minimum Off-time System\" to place a current transformer in a circuit which controls the switching of a compressor. The current transformer, in turn, provides a control signal to a minimum off-time circuit. The purpose of the delay effected by the off-time circuit is to prevent start-up under heavy load. Like the system disclosed in Harris, supra, the system is analog in nature and is responsive to analog signals.\nIt is known from U.S. Pat. No. 3,636,369 to Donald G. Harter, granted on Jan. 18, 1972 and entitled \"Refrigerant Compressor Control-relay to Control Two Time Delays\" to provide a refrigerant compressor with a relay arrangement which controls two time delays. One time delay keeps the compressor deenergized at least for a predetermined period after each stop action. The other time delay keeps the compressor energized for at least a predetermined period after each start action.\nIt is known from U.S. Pat. No. 4,033,738 to Carl R. Metrola et al., granted on July 5, 1977 and entitled \"Heat Pump System with Multi-stage Centrifugal Compressors\" to arrange multistage compressors in series. The start of the second compressor is delayed for a sufficiently long period to enable the first compressor to reach its design speed.\nThe use of a microprocessor and timed solid state logic circuitry to control automatically sequencing of compressors in a refrigeration system is disclosed in U.S. Pat. No. 4,152,902 granted May 8, 1979 to Lawrence E. Lush and entitled \"Control for Refrigerator Compressors\"."}
{"text": "This invention relates to electronic music systems, and more particularly relates to a method permitting interactive performance of music generated by an electronic music device. This invention is more specifically directed to synthesizer or computer-generated music, especially automatic or semiautomatic digital generation of music by algorithm (i.e., by computer program).\nIn the recent past, there have been proposed music generating systems, to be comprised of a digital computer and a music synthesizer coupled thereto. In performing typical such systems, the generated music is determined entirely by the user of the system, playing the role of performer or composer. The user first determines the nature of the sounds of the system produces by manipulating a plurality of controls, each associated with one or more parameters of the sound. Once the sounds are determined, the user performs music with the system in the manner of a traditional musical instrument, usually by using a piano-type keyboard.\nA major problem with the traditional approach to music as applied in the above-mentioned systems, is that it requires a considerable technical knowledge of sounds that are produced and varied electronically. Another problem is that such systems produce each sound only in response to external stimuli (i.e., acts performed by the user of the system), thereby limiting the complexity of the system's output to what the user is capable of performing. Still another problem is that the relationship between the system and user is limited to the type of functioning typical of a traditional musical instrument, so that the user can relate to the system only as a performer relates to his or her instrument. A further problem is that the peformance device employed by the user is normally a fixed part of the system, and is not interchangeable with other peformance devices.\nPrevious systems have not automatically generated sounds, music, or performance information, while allowing a performer to interact with and influence the course of the music. No previous system designed for performance could be used effectively by a performer or user not having previously learned skills, such as those required to play a keyboard instrument."}
{"text": "The instant invention relates to a modular device for transporting, aligning and stacking envelopes, and more particularly to such a device which can be combined with other like devices to provide any desired length of transport path and any desired number of stacking bins.\nIn high volume mailing applications, a plurality of documents are fed from document feeding devices to a location where they are inserted into an envelope. The envelopes are then transported seriatim from the document inserting position to another location for further processing, which typically includes the printing of postage on the envelopes. In many applications, it is desirable to separate the output of the inserter by diverting certain envelopes from the mainstream, or by directing certain envelopes, such as all those having the same zip code, to a designated collecting bin. When the objective is to print postage on the envelopes downstream of the inserting apparatus, it is necessary not only to transport the envelopes from the inserting apparatus to the postage meter or other printing device employed, but also to align the envelopes so that the printing occurs uniformly at the proper location on the envelopes.\nThere are available today various modules to perform each of the processes described hereinabove, i.e. transporting, aligning, diverting and stacking. These modules are situated at an appropriate location in the envelope path and perform their function as required. However, there is no single module available today which can be used to perform all of these tasks. Such a module obviously would provide significant cost savings to both the producer of the module and to the end user who would be able to simply purchase a plurality of similar modules and then deploy them as he saw fit.\nAccordingly, the instant invention comprises a single module to be used in the envelope transport path which can be set up in a variety of formats and modified in order to perform the functions of transporting the envelopes, aligning the envelopes, stacking the envelopes and deflecting the envelopes out of the envelope path, all at a speed which does not slow down the inserting apparatus or the postage meter or other printing apparatus used to print the postage on the envelope."}
{"text": "1. Field of the Invention\nThe present invention relates to storage managers that provide data storage services to software applications. More particularly, the invention concerns the provision of filtering functions such as encryption, compression and other data conversions as part of storage manager operations.\n2. Description of the Prior Art\nBy way of background, a storage manager is a system that acts as an intermediary between a software application (such as a backup/restore program or a web server) and a data storage resource (such as a tape drive, a disk drive, a storage subsystem, etc.). The storage manager, which could be integrated with the application program or implemented separately therefrom, provides an interface that accepts objects for storage and subsequently retrieves the objects upon request. Applications for which a storage manager has been used include the management of backup images of database installations, enterprise application data, individual workstations, web content, etc.\nThere is often a need for a storage manager to filter the data being written to or read from physical storage devices by compressing or encrypting the data. Existing storage managers that provide support for compression and/or encryption do so in one of two ways. Most commonly, such filtering is provided by algorithms that are embedded in the storage manager product itself. Less commonly, such filtering is supported by providing a programming hook that gives a storage manager user the option of writing their own algorithm(s). With this option, the user is also required to re-implement much of the functionality of the storage manager on their own.\nDrawbacks of the first approach include:                The user is limited to the compression and/or encryption algorithms that are built into the storage manager product.        Some products support encryption but not compression and vice versa.        Some products support only weak encryption or poor compression.        The storage manager vendor may charge customers extra to enable the compression and/or encryption algorithms that are built in.        If a built-in algorithm is found to have a security flaw or a crippling bug, a customer cannot immediately swap in a different off-the-shelf algorithm to avoid exposure to the risk.        Storage manager customers must wait for the vendor to update the embedded algorithms with the latest technology when better algorithms are invented, even though the new technology may already exist in stand-alone off-the-shelf programs.        A vendor may not implement a particular compression or encryption algorithm that a customer desires.        \nDrawbacks of the second approach include:                The storage manager programming hook places a burden on the customer to re-implement much of the functionality the storage manager otherwise provides. The user must typically write a program that can accept objects for storage, track the location of these objects, write and read them to/from physical storage devices, and retrieve them upon request based on whatever query protocol the storage manager requires, as well as write in the desired compression and/or encryption algorithms. In this solution, the storage manager essentially delegates all work to the user and does not provide any functionality of its own. The storage manager mostly acts as a hollow shell or “stub” that forwards all storage and retrieval requests to the user-written external program for handling. The storage manager itself merely assembles and disassembles buffers of information that pass between it and the application that is calling it, and provides stubs for the interface APIs (Application Programming Interfaces) but delegates most of the work to the user's program.        This approach provides very little support for compression and encryption. There is the programming hook but the customers are required to create the needed support at great additional expense and effort to themselves.        A customer who uses the programming hook but does not sufficiently test and debug their external program may find that their data has been corrupted by their own custom program, or that bugs in the program prevent the retrieval of storage objects at a critical time, such as when they need them to restore a down system.        If the event described in the preceding paragraph occurs, the storage manager vendor may find itself exposed to liability for the customer's own programming mistakes.        \nAccordingly, a need exists for a storage manager filtering technique that overcomes the foregoing disadvantages. What is required is a solution that allows storage manager filters to be easily implemented without having to redesign the storage manager or duplicate its functionality in a custom program. It would be further desirable to provide the capability of implementing new and different filters. At present, the most common needs for storage manager filtering are compression and encryption. However, it is submitted that the possibilities are broader, and it may be advantageous in some circumstances to provide other data conversions, such as converting between English and metric units, or between different code pages or character sets like ASCII (American Standard Code for Information Interchange), EBCDIC (Extended Binary Coded Decimal Interchange Code), and Unicode. By way of example, this capability would be useful if backup data was generated by a first system in a first character format (e.g., a mainframe computer using EBCDIC character) and the data needed to be restored to a second system that used a second character format (e.g., a workstation using ASCII character encoding). Another area where storage manager filtering could be used is the generation of audit trails. Such a filter could be used to inspect the data being stored or retrieved and generate audit information for management purposes."}
{"text": "Aircraft regulations require designers to consider flight conditions that contribute to ice formation and ice accumulation on critical portions of the aircraft. Leading edges of wings and engine nacelles can be particularly susceptible to ice formation, and require active ice protection in many aircraft designs. With respect to aircraft engines, accumulated ice can break away from the lip of the nacelle inlet, and enter the engine. Ice entering an engine can damage an engine's fan blades, or other critical engine components. Ice formation and accumulation on a nacelle inlet lip also can restrict airflow, thereby hindering engine performance. Accordingly, ice protection systems are needed for wing leading edges and engine nacelles in general, and for nacelle inlet lips in particular.\nVarious systems and methods are known for minimizing and eliminating ice accumulations on critical surfaces of aircraft. For example, the airport crews commonly spray an ethylene glycol de-icing solution on accumulated ice on aircraft wings while the aircraft are on the ground preparing for departure. Alternatively, some aircraft are equipped with pneumatically actuated bladders along leading edge surfaces of their wings that can be periodically inflated to shed accumulated ice. Many jet aircraft direct hot gases from their engine compressors onto the wing or nacelle inlet leading edges to melt accreted ice. Though such hot gas systems can generally be effective, such systems are not available for aircraft that do not have jet engines, or aircraft that do not have sufficient hot air capacity for such purposes.\nAnother method of preventing and/or eliminating ice from aircraft leading edges employs resistance-heating elements positioned along an aircraft's leading edges. Such electrothermal systems use various types of electric heating elements that are affixed on a surface structure of an aircraft. For example, the heating elements may include metallic electrodes arranged in a serpentine pattern, and affixed to a substrate that is attached to a surface structure of an aircraft. Other similar systems use ribbons or sheets of electrically conductive material as heating elements. Such systems commonly include heating elements that are intermittently spaced along an aircraft surface in a manner such that individual heaters can be selectively energized. Because most aircraft have limited available electrical power, the individual heating elements or sets of heating elements can be sequentially energized to conserve the amount of power consumed at any one time during a heating cycle.\nAn airplane's airframe and engines produce varying amounts of audible noise during takeoffs and landings. For example, an aircraft's engines typically operate at or near maximum thrust as the aircraft departs from or approaches an airport. Aircraft engine fan noise can be at least partially suppressed at the engine nacelle inlet by a noise absorbing liner. Such liners are provided inside of and proximate to the nacelle inlet. These liners can convert acoustic energy into heat, and typically consist of a porous skin supported by an open-cell honeycomb matrix. The open-cell matrix provides separation between the porous skin and a non-perforated backskin. Some have postulated that the partially open cells of the liner create a Helmholtz resonant effect that absorbs sonic energy, and thereby effectively suppresses at least a portion of the generated engine noise. Government regulators often mandate aircraft engines with reduced noise signatures, and as a result, aircraft manufacturers, airline companies, and airport communities frequently demand such engines on aircraft.\nThough electrothermal systems can be effective in preventing ice formation or shedding ice from various sensitive areas of aircraft, such systems generally do not provide for noise attenuation, such as is beneficial at the lip of an engine nacelle inlet. Conversely, prior art noise attenuation systems for aircraft generally do not provide ice protection for leading edges of the aircraft. Accordingly, there is a need for an electrothermal ice protection system for the leading edges of aircraft that also includes noise attenuation capability. In particular, there is a need for an electrothermal ice protection apparatus for a nacelle inlet noselip that is capable of attenuating at least some engine fan noise."}
{"text": "Since its introduction in 1975, the well-known Kohler and Milstein technique (Nature 256:495, 1975) for the production of mouse hybridoma cells has made it possible to produce large quantities of mouse antigen-specific monoclonal antibodies that are useful in a number of investigative, diagnostic and therapeutic applications. The mouse hybridoma cells, which are initially produced by the fusion of antibody-producing cells (B-lymphocyte cells, hereinafter referred to as B-cells) with malignant, transformed B-cells (in vivo transformed, myeloma cells from mice afflicted with myeloma or plasmacytoma) are capable of producing large quantities of monoclonal antibodies with predetermined specificities.\nUsing the Kohler and Milstein technique, a B-cell and a plasmacytoma cell are fused using, for instance, polyethylene glycol, lysolecithin or Sendai virus as the cell-fusing agents. A selectable marker must be present in the fused cells to enable them to be selected from parent cells and other non-hybridoma cells. As an example, the plasmacytoma fusion partner is generally deficient in an enzyme, (for instance, hypoxanthine-guanosyl phosphoribotransferase (HGPRT)) that is necessary for growth of the fused cell in certain media (hypoxanthine-, aminoprotein-, and thymidine- containing medium or HAT medium). This enzyme deficiency enables the resultant hybrids to be selected for their ability to grow in such media. This insures that only B-cell:plasmacytoma cell hybrids are recovered since neither parental cells (nor hybrids comprising B-cell: B-cell and plasmacytoma: plasmacytoma cell) can survive in selective media.\nMurine antibodies produced with the Kohler and Milstein technique are generally unsuitable for administration to human subjects as in-vivo therapeutic agents, e.g., to provide passive immunity to an infectious agent. The extension of the Kohler and Milstein hybridoma technology to the production of human monoclonal antibodies has been limited, largely due to: (1) the lack of good human plasmacytoma cells for use as fusion partners; (2) the low frequency of cell fusion events (\"fusion efficiency\"); and (3) the relative scarcity of B-cells circulating in human blood and producing specific antibodies against antigens of interest (and the inherent difficulties in isolating such cells). These factors make it difficult to obtain hybridoma cell lines secreting human monoclonal antibodies of a predetermined specificity.\nCasali et al. (Science 234:476-479, 1986) disclosed a method which represents some progress toward making human monoclonal antibody-producing cells. Normal B-cells obtained from peripheral human blood were pre-selected for their specificity to a given antigen using Fluorescence-Activated Cell Sorting (FACS). Positively selected clones were then established as lymphoblastoid cells in vitro by infecting such cells with Epstein-Barr virus (EBV). The EBV infected cells produced antigen-specific human monoclonal antibodies. However, the method of Casali et al. has several significant drawbacks which impair its usefulness: (1) the amount of monoclonal antibodies produced by the Casali et al. cells is relatively low, and (2) the antibody producing cells are relatively unstable and some clones stop antibody production prematurely. In addition maintenance of the antigen-specific antibody production requires repeated cloning of the cells, a time-consuming and inefficient procedure given the low clonogenic (i.e. growth) properties of the resultant lymphoblastoid or lymphoblastoid cell lines (LCL); (3) large-scale production and purification of the monoclonal antibodies is inefficient in view of the long doubling time and high serum requirements of the LCL; and (4) the LCL produced by this process cannot be grown as tumors in animals. Such tumor cell growth permits the amplification and purification of antibodies from ascitic fluids, an efficient method for large scale antibody production that is widely used in making murine monoclonal antibodies. Finally, the Casali et al. method does not dispense with the requirement for identifying a human B-cell specific to a certain antigen.\nCopending U.S. patent application Ser. No. 041,803 (allowed) filed Apr. 23, 1987 of Riccardo Dalla-Favera discloses a method for the production of human monoclonal antibody-producing cells. Specific B lymphocytes are selected using the method of Casali et al. (supra), infected with Epstein Barr virus (EBV) and transfected with activated c-myc DNA sequences. The resultant cells are tumorigenic (i.e. can grow in semisolid medium and animals such as rats or mice) and clonogenic and produce monoclonal antibodies of a predetermined specificity. However, it was found that these cells still produce relatively low amounts of antibody because the transfected lymphoblastoid cells had not undergone differentiation.\nCurrently there is no convenient and reliable system available for the production of human monoclonal antibodies wherein the monoclonal antibody-producing cells are stable, highly malignant and which can be readily manipulated to produce high antibody titers.\nIt, is therefore an object of the present invention to provide a method for the production of tumorigenic human cells that are capable of producing human monoclonal antibodies.\nA further object of the present invention is to provide a transformed lymphoblastoid cell that is useful as a fusion partner in the production of human monoclonal antibodies.\nAnother object of the present invention is to provide a transformed lymphoblastoid cell that demonstrates high level proliferative, differentiation and antibody production properties.\nAnother object of the present invention is to produce a new human cell line comprising human B-cells infected with Epstein-Barr virus and which have at least one exogenous activated K-, N- or H-ras oncogene DNA sequences."}
{"text": "Targeting disease-causing gene sequences was first suggested more than thirty years ago (Belikova et al., Tet. Lett., 1967, 37, 3557-3562), and antisense activity was demonstrated in cell culture more than a decade later (Zamecnik et al., Proc. Natl. Acad. Sci. U.S.A., 1978, 75, 280-284). One advantage of antisense technology in the treatment of a disease or condition that stems from a disease-causing gene is that it is a direct genetic approach that has the ability to modulate (increase or decrease) the expression of specific disease-causing genes. Another advantage is that validation of a therapeutic target using antisense compounds results in direct and immediate discovery of the drug candidate; the antisense compound is the potential therapeutic agent.\nGenerally, the principle behind antisense technology is that an antisense compound hybridizes to a target nucleic acid and modulates gene expression activities or function, such as transcription and/or translation. The modulation of gene expression can be achieved by, for example, target degradation or occupancy-based inhibition. An example of modulation of RNA target function by degradation is RNase H-based degradation of the target RNA upon hybridization with a DNA-like antisense compound. Another example of modulation of gene expression by target degradation is RNA interference (RNAi). RNAi generally refers to antisense-mediated gene silencing involving the introduction of dsRNA leading to the sequence-specific reduction of targeted endogenous mRNA levels.\nAn additional example of modulation of RNA target function by an occupancy-based mechanism is modulation of microRNA function. MicroRNAs are small non-coding RNAs that regulate the expression of protein-coding RNAs. The binding of an antisense compound to a microRNA prevents the microRNA from binding to its messenger RNA target, and thus interferes with the function of the microRNA. Regardless of the specific mechanism, this sequence-specificity makes antisense compounds extremely attractive as tools for target validation and gene functionalization, as well as therapeutics to selectively modulate the expression of genes involved in the pathogenesis of malignancies and other diseases.\nAntisense technology is an effective means for reducing the expression of one or more specific gene products and can therefore prove to be uniquely useful in a number of therapeutic, diagnostic, and research applications. Chemically modified nucleosides are routinely incorporated into antisense compounds to enhance one or more properties, such as nuclease resistance, pharmacokinetics or affinity for a target RNA. In 1998, the antisense compound, Vitravene® (fomivirsen; developed by Isis Pharmaceuticals Inc., Carlsbad, Calif.) was the first antisense drug to achieve marketing clearance from the U.S. Food and Drug Administration (FDA), and is currently a treatment of cytomegalovirus (CMV)-induced retinitis in AIDS patients.\nNew chemical modifications have improved the potency and efficacy of antisense compounds, uncovering the potential for oral delivery as well as enhancing subcutaneous administration, decreasing potential for side effects, and leading to improvements in patient convenience. Chemical modifications increasing potency of antisense compounds allow administration of lower doses, which reduces the potential for toxicity, as well as decreasing overall cost of therapy. Modifications increasing the resistance to degradation result in slower clearance from the body, allowing for less frequent dosing. Different types of chemical modifications can be combined in one compound to further optimize the compound's efficacy.\nVarious antiviral monomers based on the bicyclo[3.1.0]hexane pseudo-sugar analog scaffold have been reported (see PCT International Application WO 2006/128159, published on Nov. 30, 2006; PCT International Application WO 2006/091905, published on Aug. 31, 2006; PCT International Application WO 01/51490, published on Jul. 19, 2001; PCT International Application WO 95/08541, published on Mar. 30, 1995; PCT International Application WO 02/08204, published on Jan. 31, 2002; and Kim et al., J. Med. Chem., 2003, 46, 4974-4987).\nsiRNAs with one or two ribo-like north bicyclo[3.1.0]hexane pseudosugars (2′-OH) have been prepared (see Terrazas et al., Organic Letters, 2011, 13(11), 2888-2891). The Tms of the resulting oligos was lowered by addition of the modified pseudo-sugar analogs (−1.6° C./modification). In vitro studies using the siRNA with one or two modifications compared to wild type guide strand showed that one incorporation had comparable results to wild type and two modifications was less active.\nOligonucleotides have been prepared with one or two ribo-like north bicyclo[3.1.0]hexane pseudosugars (2′-H) (see Maier et al., Nucleic Acids Research, 2004, 32(12), 3642-3650)."}
{"text": "My invention relates to a signal restoration circuit, and particularly to a signal restoration circuit for accurately restoring distorted binary digits or bits to their original form.\nBinary data is used extensively for signalling and information purposes. While such binary data has many advantages, it is subject to distortion or error, particularly when transmitted over a medium such as a radio system. Consequently, there is a need for an arrangement to restore the binary digits or bits to their original form.\nAccordingly, a primary object of my invention is to provide a binary digit or bit restoration circuit.\nAnother object of my invention is to provide a new and improved circuit that accurately restores received bits to their original form despite severe distortion of the bits by the transmission medium.\nAnother object of my invention is to provide a new and improved bit restoration circuit that can be implemented with digital circuits.\nIn addition to being subject to distortion during transmission, binary digits or bits are liable to be improperly interpreted or detected at their leading and trailing edges when the relevant circuits are responding to those edges.\nAccordingly, another object of my invention is to provide a new and improved bit restoration circuit which, in determining the logic of received bits, omits consideration of the bits in the vicinity of their leading and trailing edges."}
{"text": "1. FIELD OF THE INVENTION\nThe present invention relates to an automatic layout apparatus and method of automatically selecting and laying out layout candidate elements in an analog LSI.\n2. DESCRIPTION OF THE PRIOR ART\nIn an analog LSI layout design, a differential pair or current-mirror-connected bipolar transistors are laid out close to each other in consideration of the layout direction so as to satisfy the relative precision of element characteristics in order to satisfy the electrical characteristics of a circuit.\nThe feature in this case is common connection of emitters or bases. When a measure against crosstalk must be taken depending on a circuit arrangement, a layout in which lines need not be crossed is performed in a line design upon a layout design. When a measure against a heat source is necessary, a power transistor serving as the heat source and a signal system transistor having electrical characteristics changed by the heat must be laid out apart from each other.\nSince the analog LSI layout design has many limitations including these items on the electrical characteristics of the circuit, it has manually been performed. However, as the analog LSI becomes larger in scale year by year, the design becomes more difficult to perform by the method of manually selecting layout candidate elements one by one such that the development period becomes longer. For this reason, the design period is tried to be shortened by, e.g., using the layout candidate element automatic selection function of an automatic layout/wiring apparatus used for a MOS or the like. Although layout candidate elements can be automatically selected by using the automatic layout/wiring apparatus, automatic selection is independent of the electrical characteristics and the intention of the layout designer because the automatic selection method is selection based on the element number order and selection based on the description order in a circuit information file. Accordingly, the layout designer must manually select elements again, so the design efficiency cannot be improved.\nAn example of a conventional layout flow in the use of the automatic layout/wiring apparatus will be described with reference to a flow chart shown in FIG. 1 and a circuit example shown in FIG. 6.\nIn the first step, connection information of a circuit shown in FIG. 6 is input (F-1). In the second step, an arbitrary element is manually selected and laid out (F-2). In this case, a transistor Q1 is selected and laid out. In the third step, a transistor Q2 is selected as a next layout candidate element in a tool of selecting elements in the element number order (F-3). Finally, in the fourth step, whether the transistor Q2 selected in the third step is proper is checked (F-4).\nSince a transistor Q6 constituting a differential pair with the transistor Q1 is wanted to be laid out in an original analog LSI layout, the layout designer selects either of\n1 selection is canceled, and the flow returns to the second step to manually select the transistor Q6; and PA1 2 the transistor Q2 is tentatively laid out in the fifth step (F-5), and the flow returns to the third step to automatically, newly select a next layout candidate element. When tentative layout is performed, manual selection and re-layout must be performed at a proper time. In this manner, the processing from the second step to the fifth step is repeatedly executed until all elements are laid out.\nAs described above, the conventional operation is very cumbersome such that an automatically selected element is canceled to manually repeat selection/layout, or re-layout is performed at a proper time upon temporary tentative layout. Various automatic layout methods using an automatic layout/wiring apparatus are examined. In these methods, constraint items are registered in a database by, e.g., designating a relative layout definition or the like serving as a constraint item based on electrical characteristics unique to an analog circuit for a group in units of blocks, and automatic layout is performed on the basis of the database. This database has no versatility, and must be prepared for each product, resulting in a low efficiency.\nThe first problem of the prior art is therefore a low design efficiency. In the conventional manual layout method, layout candidate elements must be selected one by one, or when the layout candidate element automatic selection function of the automatic layout/wiring apparatus is used, layout candidate elements are automatically selected from elements which are not laid out yet. However, an element intended by the layout designer in consideration of a layout that satisfies the electrical characteristics such as the relative or absolute precision is not always selected as a layout candidate because the selection method is executed in the element number order or the description order in a connection information file. For this reason, an element must be manually selected again, or an optimal layout must be performed by temporarily performing tentative layout and selecting an element again at proper time.\nThe second problem of the prior art is that a database (limitation item) used in performing automatic layout by the conventional automatic layout/wiring apparatus has no versatility, and must be prepared for each product because the database is formed mainly using the element numbers or the element numbers of groups. That is, since the element numbers of respective products are different from each other, the database must be prepared for each product. Even in the same product, the database must be constructed again upon changing the circuit."}
{"text": "Signal processing architectures are one of the main foundational components of the modern digital age. As is common in ordinary desktop or mobile computer applications, users are given a plurality of multimedia choices when viewing, listening, and/or interacting with data that has been processed by such systems. Before users actually utilize such data in a respective application, however, analog information is typically sampled and captured in real time via an analog-to-digital converter and processed via a Fast Fourier Transform (FFT) and/or other signal processing techniques. Sampled data is often stored in a database whereby subsequent signal processing and/or data manipulation is performed thereon. After the data has been stored, a plurality of database algorithms or techniques may be employed to retrieve such data and are described below. Unfortunately, the form of data storage such as via a floating-point format is not very conducive to efficient processing and retrieval of the data. Moreover, noise that may be present in any given sample of data may cause significant problems when determining if another previously stored and/or related sample can be located in the database. For example, if a recently captured data sample were sent to a database of stored samples that are potentially related to the captured data, and the recently captured data was taken in a noisy environment, it may be substantially difficult (or not possible) to determine if the noisy sample matches or relates to any of the previously stored samples in the database (e.g., require large amounts of processing bandwidth to determine a match, if any).\nAs noted above, many database techniques have evolved to locate and retrieve previously stored data such as can be provided by various tree lookup procedures. For example, there are many variants of tree lookup processes that attempt to speed-up basic nearest neighbor determinations. One of the earliest known is the k-d tree, which is a binary tree wherein the data is split, according to the value of a particular component, such that roughly half of the data falls on either side of the split, whereby the particular component is selected to maximize the variance of the data in a direction perpendicular to a corresponding hyperplane. In a test phase, a rectangle containing a test point is located by descending the tree, wherein backtracking (e.g., process of retracing a search path) is performed if the closest training point in an associated hyperectangle is such that points in adjacent rectangles may be closer. It is believed that k-d trees are somewhat limited to applications having lower dimensional structures (e.g., about 10 dimensions). In addition, the k-d tree has the property that rejection (of a point that falls farther than a threshold away from all other points in the database) can be as computationally expensive as finding the nearest neighbor.\nMore recently, a variety of trees—an R-tree, an R* variant, and for example S—S trees have been proposed. In these trees, processed nodes correspond to regions in space into which the data falls, so if a test point falls in a node, the other points in that node are known or assumed to be close to the test point. However, this does not obviate the need for backtracking, but facilitates making an early rejection possible—a property that k-d trees do not have. In R-trees, the nodes are populated by rectangles. R-trees are a variant that tend to minimize the area, margin and overlap of the rectangles (whereby the ‘margin’ of a rectangle may be defined as the sum of the lengths its sides), which generally results in faster lookup, and also introduces ‘forced reinsertion’, for providing a more balanced tree.\nThe S—S (similarity search) tree approach may even out-perform R-trees on high dimensional data. In this approach, leaves of the tree correspond to ellipsoids, in which a center and radius are defined by the data enclosed (generally, the principal axes of the ellipsoid are selected beforehand, and represent the relative importance of different dimensions). The center of the ellipsoid is thus, the centroid of the data, wherein the radius is selected to enclose the data. Again, forced reinsertion is employed to balance the tree. Other approaches have focused on how approximate matching (that is, given a query q and some set of points P, find a point pεP such that ∀p′εP,d(p,q)<(1+e)d(p′,q), for some small e, wherein d(p,q) is a distance measure between p and q) can yield more optimal bounds on preprocessing and lookup times than exact matches provide, however, the lookup time scales as (1/e)d, wherein d is the dimension of the space which may cause an impractical computational expense for many applications that employ higher dimensional data sets."}
{"text": "The invention concerns cosmetic compositions for the treatment of hair or skin, having a content of new, macromolecular compounds derived from chitosan, which are employed in a suitable cosmetic foundation.\nThe invention further concerns new N-hydroxypropyl-chitosans, as well as processes for the production thereof.\nIt is already known to employ cationic polymers, in particular polymers which display quaternary ammonium groups, as conditioning agent in cosmetic compositions, particularly for the treatment of hair. Based upon a reciprocal action between their ammonium groups and the anionic groups of the hair, the cationic polymers possess a great affinity towards keratin fibers.\nIt has been determined that the employment of such cation-active polymers in such cosmetic compositions provides numerous advantages: the disentanglement of the hair, as well as its treatment, are facilitated, and, moreover, the hair is provided with elasticity and lustrous effect. On account of their affinity towards keratin, however, these polymers tend to accumulate in the hair upon repeated use, so that the hair becomes heavier, which is undesirable as a final effect.\nMoreover, synthetic polymers provide problems on account of the physiological activity of possibly present trace monomers, which are removable from the polymers only with difficulty.\nIt has already been attempted to eliminate the above-mentioned disadvantages by emplying in such cosmetic compositions the water-soluble salts of chitosan, i.e. polyglucosamines producable from chitin by means of entacetylation. In this connection, reference is made to European Patent 0 002 506, as well as German Pat. No. 26 27 419.\nIn the same manner as with the plurality of cation-active polymers having quaternary groupings, chitosan likewise frequently provides the disadvantage that it is not too compatible with the anion-active surface-active agents which in customary manner find use in cosmetic compositions for the treatment of hair, particularly in shampoos. It is therefore necessary to apply the chitosan for penetration in separate treatments, namely before and/or after the shampooing.\nIn addition, the chitosan displays, in neutral and alkaline medium, near insolubility, so that its use, for example, in alkaline permanent shaping compositions or hair dyeing compositions, is not possible.\nBy means of employment of glycidyl chitosans instead of chitosan salts according to DE-OS 32 No. 23 423, the above-mentioned disadvantages can be avoided. The reaction of chitosan with glycide is, however, very cost-intensive, since glycide is a more expensive raw material, not produced on a large scale."}
{"text": "The term “prion” was coined in 1982 by Stanley Prusiner when describing his findings on the causative agents of the transmissible spongiform encephalopathies: proteinaceous infectious particles that lack any nucleic acid. These agents are truly unprecedented in medical science as infectious pathogens that cause fatal neurodegenerative disorders through an entirely novel mechanism. While prions may present as infectious, genetic, or sporadic disorders, they all develop as a direct result of a biochemical modification to the prion protein (PrP) that is a normal constituent of all mammalian cells (Prusiner, S. B., Proc. Natl. A cad. Sci. USA 9513363-13383(1998)).\nThe earliest studies of scrapie pathogenesis demonstrated that the disease could be directly transmitted from one animal to another. The scrapie agent was originally believed to be a virus, but it has, unlike known animal or any other kind of viruses, many unique characteristics such as the extraordinarily long incubation period to disease, the noninflammatory degenerative abnormalities that developed in the brain, and the lack of any demonstrable virion particles by classical virological techniques. The long incubation periods of the scrapie agent sets them apart from most viral infections. The complete absence of a detectable immune response is puzzling but this may now be explained by the fact that the agent may be a modified host protein. PrP is a host-specific protein, encoded by a single exon of a unique host gene. PrP is the product of highly conserved gene found in diverse organisms, and is a membrane bound protein thought to have an important, but yet unknown function.\nBrains of scrapie-infected hamsters contain two forms of PrP: the cellular PrP (PrPC) and the scrapie PrP (PrPSc) isoforms. Both proteins have a mass of 33-35 KD but they have different physical properties. PrPC is anchored to the cell surface and can be solubilized with ionic detergents as well as being susceptible to proteolytic agents. In contrast, PrPSc cannot be solubilized by ionic detergents and looses only an amino-terminal peptide to proteolytic agents to yield a protein of mass 27-30 KD called PrP 27-30 (Prusiner et al. Cell 38:127-140(1994)).\nPrP 27-30 is the major constituent of the pathognomonic amyloid plaques that are found in the brains of many hosts with spongiform encephalopathies. The quantity of this novel protein correlated with the titer of prion infectivity in brain. Moreover, PrP 27-30 was absent from uninfected brain, and it was found that various procedures that denatured, hydrolysed, or modified PrP 27-30 also inactivated prion infectivity.\nNo differences in the primary structure (i.e. amino acid sequence) of PrPC and PrPSc have been detected, nor have any differences been found between PrP genes or mRNAs from normal and infected brains with respect to structure or copy number. The physical differences such as three-dimensional configuration between the two proteins are therefore attributed to post-translational chemical modification. In general, during the refolding of PrPC into PrPSc, some of the normal α-helical protein structure is partially converted into β-sheet.\nTo describe the nature of scrapie agent, two hypotheses were proposed: 1) a “protein only” hypothesis, in which the prion particle is devoid completely of nucleic acid; and 2) a “nucleoprotein or virino” hypothesis, in which the prion consists of a small nucleic acid and host-encoded protein. Sparrer et al.'s experiments suggest that protein only hypothesis is correct by using a yeast prion-like system (Sparrer H. E. et al. Science 289, 595 (2000)).\nThe recent pathogenesis studies of bovine spongiform encephalopathy have shown that experimentally infected cattle can show prion infectivity in the ileum (small intestine) in advance of their neurologic disease (Collee et al., Lancet 349:636-640 (1997)). Epidemiologic data now support an oral route of transmission in a number of animal prion disease outbreaks, although how sporadic prion diseases, such as Creutzfeldt-Jakob disease in humans, develop still remains unknown. Nevertheless, the fact that brain tissue from an affected host can transmit disease to an unaffected recipient (particularly if such material is inoculated directly into the brain of that recipient) now stands as one of the defining characteristics of all prion diseases. In addition, it has become clear that scrapie can induce disease in rodents following either a peripheral (subcutaneous, intraperitoneal, and oral) or an intracerebral inoculation.\nIt was reported that prion neuroinvasion requires B lymphocytes (Klein M A et al. Nature 390:687-690 (1997)). Almost paradoxically, normal PrP expression is not required for B cells to transmit disease to the brain, suggesting either that other cell type(s) whose maturation depends on B cells or their products (such as follicular dendritic cells) may promote neuroinvasion, or that B cells carry prions to the nervous system in a PrP-independent manner (Klein M A et al. Nature Med 4:1429-1433 (1998)).\nIn initial experiments, it was demonstrated that a prion disease in one species could be transferred to another. However, subsequent attempts at cross-species transmission were inconsistent. Recently, it has become clear that the successful passage of prions between species is almost always characterized by a prolonged incubation period during the first passage in the new host.\nThis time delay is often referred to as the prion “species barrier.” However, on the next passage into a homologous host, the incubation period shortens and remains remarkably constant for all subsequent passages in that species. The species barrier occurs because new prions synthesized de novo in an experimentally inoculated host are generated from, and therefore reflect the protein sequence of, the host PrP and not that of the PrPSc molecules in the inoculum (Bockman J M et al. Ann Neurol. 21:589-595 (1987)). Thus, the prion donor is the last mammal in which the prion was passaged and its PrP sequence represents the “species” of the prion.\nThe prion species is differs from the prion “strain” whose information appears to be enciphered in the conformation of the nascent PrPSc. Both the species and strain influence the ability of a given prion to cause symptomatic disease in a heterologous host. The species barrier concept is of practical importance in assessing the risk that humans may develop a prion disease after consuming scrapie-infected lamb or bovine spongiform encephalopathy-infected beef.\nIn yeast, two notable prion-like determinants [URE3] and [PSI], have been described (Reed B. Wickner, Science 264, 566-569 (1994); Wickner et al., J. Biol. Chem 274(2), 555-558 (1999)). Interestingly, different strains of yeast prions have been identified. Conversion to the prion-like [PSI] state in yeast requires the molecular chaperone Hsp 104; however, no homolog of Hsp 104 has been found in mammals (Patino, M. M et al., Science 273, 622-626(1996)). The NH2 terminal prion domains of Ure2p and Sup35 that are responsible for the [URE3] and [PSI] phenotypes in yeast have been identified. In contrast to PrP, which is a GPI-anchored membrane protein, both of Ure2p and Sup35 are cytosolic proteins (Reed B. Wickner,. Proc. Natl. Acad. Sci. USA 94, 10012-10014(1997)).\nThere have been efforts to provide treatment of prion diseases. For example, Tomiyama et al. disclose that antibiotic, rifampicin and its derivatives, which possess a naphthohydroquinone or naphthoquinone structure, inhibited Aβ1-40 aggregation and neurotoxicity in a concentration-dependent manner. Hydroquinone, p-benzoquinone and 1,4-dihydroxynaphthalene also inhibit Aβ1-40 aggregation and neurotoxicity at comparable molar concentrations to rifampicin (Tomiyama et al. JBC 271 (12): 6839-6844 (1996)).\nCaspi et al. showed anions such as Congo red(CR) reduce the accumulation of PrPSc in a neuroblastoma cell line permanently infected with prions as well as to delay disease onset in rodents when administered prophylactically.(Sigal Caspi et al., The Anti-prion Activity of Congo Red, The Journal of Biological Chemistry, 273(6), 3484-3489 (1998)).\nDE 4229805 discloses that toxic effects displayed by PrPSc and its peptide fragment can be blocked by antagonists of N-methyl-D-aspartate (NMDA) receptor channels, like Memantine. Flupirtine, a non-opiod analgesic drug, which is already in clinical use, was found to display in vitro a strong cytoprotective effect on neurons treated with PrPC or PrP106-126.\nWIPO International Patent Publication No. WOOO09111 discloses treatments of amyloidogenic diseases and prion diseases associated with conversion of protease sensitive PrP(PrP-sen) to protease resistant PrP(PrP-res) by administering tetrapyrrole such as phthalocyaninines, deuteroporphyrins and meso-substituted prophines.\nDE 4330388 discloses a curing or prevention of AIDS or mad cow disease by using L-tryptophan, indole, 3-indolylacetic acid or indomethacin to increase indole levels.\nU.S. Pat. No. 5,935,927 to Vitek et al. teaches a method for stimulating amyloid removal in amyloidogenic diseases using advanced glycosylation endproducts to increase the activity of scavenger cells within the body at recognizing and removing amyloid deposits from affected tissues and organs.\nU.S. Pat. No. 6,020,537 to Prusiner discloses prion protein standards for use as reference materials for prion detection and methods for the preparation of the prion protein standard. U.S. Pat. No. 5,962,669 to Prusiner describes a protein designated Prion Protein Modulator Factor (PPMF), which is an auxiliary factor in prion replication.\nU.S. Pat. No. 5,750,361 discloses a method of screening for compounds which inhibits the binding of PrPSc to a PrP peptide based on the fact that PrPSc an increased β-sheet content, a diminished aqueous solubility, and a resistance to proteolytic digestion, relative to PrPC.\nU.S. Pat. No. 5,948,763 to Soto-Jara et al. discloses peptides capable of interacting with a hydrophobic structural determinant on a protein or peptide for amyloid or amyloid-like deposit formation inhibit and structurally block the abnormal folding of proteins and peptides into amyloid or amyloid-like deposits.\nReferences in the art of treatment of prion diseases are scarce. Recently, U.S. Pat. No. 6,060,293 to Bohr et al. proposes treating prion related diseases by changing the functionality of the three-dimensional structure of proteins by applying high frequency energy having maximum frequency in the range of 0.01-100 GHz to a fluid system containing such proteins. However, it is not clear from the '293 patent how one would go about treating a mammal with such a disease.\nTherefore, a need still exists for treating prion diseases, such as CJD and mad cow disease."}
{"text": "This invention relates to an electric surgical knife device.\nA prior art electric surgical knife device supplies an incision current for incising body tissues, a coagulation current for coagulating the blood at the area of the incision and a blend current, i.e. blend of incision and coagulation currents, for both incising body tissues and coagulating the blood at the area of the incision. The incision current, the coagulation current and the blend current are supplied independently of one another by actuating respectively provided switches. An operating physician may mistakenly actuate two of these switches or all these switches at the same time. If this happens, the device may fail to supply the desired current. For example, when the physician mistakenly actuates the incision current supply switch and the coagulation current supply switch at the same time though he intends to coagulate the blood, the incision current may be supplied. In this case, the body tissue is cut without a blood coagulation procedure which will result in an enormous bleeding.\nIt is an object of this invention to provide an electric surgical knife device which helps carry out electric surgery in a safe and sound manner, if erroneously operated."}
{"text": "In recent years, attention has been given to sealing joints in moving vehicles, and more particularly to an apparatus for sealing a track joint in a track chain. For example, crawler tractors, such as a bulldozer, typically have a sprocket, an idler, a track chain and a number of track shoes attached to the track chain for propelling the tractor over the ground. During use of the crawler tractor the sprocket rotates and engages the track chain, thereby causing the track chain, along with the attached track shoes, to rotate around a path defined by the sprocket and the idler. The rotation of the track chain causes the track shoes to engage the ground, thereby propelling the crawler tractor over the ground to perform various work functions.\nTrack chains generally include a pair of parallel chains, with each parallel chain being made up of a series of entrained master links and track links. Some track chains may further include a series of pins and bushings interposed between and connected to the parallel chains. The bushings and the entrained track links cooperate to form a number of track joints which allow the necessary movement of the bushings relative to the track links during use of the track chain, for example when the track chain rotates about the sprocket and the idler. Track joints are typically equipped with a track seal assembly to keep out various corrosive and abrasive mixtures of water, dirt, sand, rock or other mineral or chemical elements to which the track chain is exposed during its use. The track seal assembly also functions to keep a lubricant within the track joint to facilitate the aforementioned relative movement of the bushings and the track links.\nA problem with track seal assemblies is to keep dirt out and keep lubrication within the track joint. Mud packing resistance should preferably be kept as high as possible. Another problem to consider is to maintain sealing capability. It is desired to maintain the different parts of a track seal assembly in place to avoid displacement of any parts and resulting loss of sealing capability. Wear life of a seal assembly may be a problem. The stability of a load ring and/or seal ring within a seal assembly may be a problem. A problem with seal assemblies may be to take up or compensate for misalignments within a seal assembly.\nOther problems to consider alone or in combination with other mentioned problems relate to grooving of the different parts of a seal assembly. Grooving may arise by various abrasive particles found in the working environment of a track chain finding its way in to a seal assembly. If such grooving arises, then not only dirt may enter a track joint, but also lubrication within the track joint may escape.\nIt is further desirable to keep costs for manufacturing and/or assembling a seal assembly within a track joint down. It may also be desirable to avoid cumbersome arrangements. It may also be desirable to provide a simple assembling process of a track joint seal assembly.\nThe present invention is directed to overcoming one or more of the problems as set forth above."}
{"text": "1. Field of the Invention\nThe present invention relates to a zoom lens and an image projection apparatus including the same, suitable for, for example, a liquid crystal projector having a long back focus and maintaining good pupil consistency with a lighting system.\n2. Description of the Related Art\nHitherto, various kinds of liquid crystal projectors (image projection apparatuses), each including a display device such as a liquid crystal display and projecting an image formed in the display device onto a screen, have been proposed.\nIn particular, the liquid crystal projector is in widespread use for a conference and a home theater, as an apparatus capable of projecting an image outputted from a personal computer or the like onto a large screen. A projection lens for use in the liquid crystal projector is requested to have, for example, the following features.\n(a) In a three-panel color liquid-crystal projector including three liquid crystal displays, light emitted from a white light source is generally separated into red, green and blue colors by a color-separation optical system and introduced into the corresponding liquid crystal displays. Three kinds of light emitted from the respective liquid crystal displays are synthesized by a color-synthesis optical system and incident on a projection lens.\nDue to its configuration, a space having a prism and the like arranged therein, for synthesizing the three kinds of color light passing through the liquid crystal displays, must be provided between the liquid crystal displays and the projection lens, thereby causing the projection lens to have a certain length of back focus. In addition, use of a reflective liquid-crystal display, i.e. use of a liquid crystal on silicon (LCOS) as the liquid crystal display causes the projection lens to have a longer back focus than upon use of a transmissive liquid crystal display.\n(b) When the angle of a light flux emitted from the liquid crystal display and incident on the color-synthesis optical system is changed, the spectral transmittance of the color-synthesis optical system is accordingly changed. As a result, a brightness of each color of the projected image is changed in accordance with the angle of view, thereby causing an image hard to be viewed. In order to reduce the influence of the angle dependency, the projection lens must be a so-called telecentric optical system in which a pupil close to the liquid crystal display (a reduction conjugate surface) lies substantially at infinity.\n(c) When pictures (images) of the three color liquid-crystal displays are synthesized and projected onto a screen, pixels of the respective colors must be overlaid one another across the overall screen so as prevent a loss of a resolution sensation caused by phenomena, for example, seeing two of a character.\nTo achieve this, a color drift (a chromatic aberration of magnification) generated in the projection lens must be satisfactorily compensated for in the visible light zone.\n(d) A distortion aberration must be satisfactorily compensated for so as to prevent the projected image from being hard to be viewed because of distortion.\n(e) The projection lens must have a small F-number (hereinafter, referred to as an Fno) and be a bright one so as to efficiently take in light emitted from the light source.\n(f) Recently, needs for a higher brightness and a high resolution of a projector are present, and at the same time, a projector having small liquid crystal panels installed therein is required to have a reduced size and weight for giving greater importance to its portability and mobility. At the same time, taking account of a use environment of such a small-sized projection apparatus, the projection lens must be a large magnification-varying zoom lens with which an image is projected onto a large screen with a shorter projection distance and the size of a projection screen is easily adjusted.\nThe liquid crystal projector-use projection lens required to have a long back focus and a variety of high optical properties as described above sometimes has an aspherical lens (a lens having an aspherical surface) incorporated therein as an effective mechanism for satisfactorily compensating various aberrations without increasing the number of lenses.\nAlso, when focusing attention on easy workability and productability, the aspherical lens should be composed of a plastic material.\nHitherto, a variety of liquid crystal projector-use projection lenses are proposed (see Japanese Patent Laid-Open Nos. 2004-004964, 2001-100100, and 2002-131636).\nWhile Japanese Patent Laid-Open No. 2004-004964 discloses an optical system aiming at a reduced size and a low cost, the system is directed to use in a transmissive liquid-crystal projector and, in addition to a limit for achieving a wider angle and a larger magnification, has an excessively short back focus when used in a reflective liquid-crystal projector.\nWhile Japanese Patent Laid-Open No. 2001-100100 discloses an optical system having a reduced size and inhibiting various aberrations generated due to its magnification of varying, in addition to a limit for achieving a larger magnification, the system has an excessively small back focus when used in the reflective liquid-crystal projector.\nWhile Japanese Patent Laid-Open No. 2002-131636 discloses an optical system having a long back focus, taking account of use for in the reflective liquid-crystal projector while keeping a wide angle of view, the system has a large Fno and includes a single focal-point lens.\nAlso, as a projection lens for use in the liquid crystal projector, US Published Application. 2001050818 discloses a six-group zoom lens configured by six lens groups as a whole, having a structure in which first to sixth lens groups respectively having negative, positive, positive, negative, positive (or negative), and positive refractive power are arranged in order from its magnification conjugate side (from its front side), and performing zooming by appropriately moving a predetermined lens groups of these lens groups.\nWith a structure in which the first and sixth lens groups are fixed and all of the inside second to fifth lens groups are moved towards the reduction conjugate side (the rear side of the zoom lens) upon zooming from its wide angle end to telephoto end, the above six-group zoom lens maintains the overall length constant at the time of zooming and serves as a telecentric zoom lens towards the reduction conjugate side while reducing the distortion aberration and the chromatic aberration upon zooming.\nOther than the above-described lenses, as a projection lens for use in the known liquid crystal projector, Japanese Patent Laid-Open No. 2001-108900 discloses a six-group zoom lens configured by six lens groups as a whole, having a structure in which first to sixth lens groups respectively having negative, positive, positive, negative, positive, and positive refractive power are arranged in order from its magnification conjugate side (from its front side), and performing zooming by appropriately moving a predetermined lens groups of these lens groups.\nWith a structure in which the first, fourth, and sixth lens groups are fixed and the inside second, third, and fifth lens groups are moved upon magnification of varying from its wide angle end to telephoto end, the six-group zoom lens maintains the overall length constant and serves as a telecentric zoom lens on the reduction conjugate side while inhibiting variations in various aberrations including the chromatic aberration during magnification of varying.\nAlso, as a zoom lens for use in the liquid crystal projector, US Published Application. 2003117716 discloses a five-group zoom lens having a structure in which, in order from its front to rear sides, five lens groups respectively having negative, positive, negative, positive, and positive refractive power, and performing zooming by moving a plurality of the lens groups of these lens groups.\nAlso, while a zoom lens disclosed in Japanese Patent Laid-Open No. 2004-138678 has an optical system having a relatively long back focus, its Fno at its wide angle end is large on the order of 2.3 or 2.4, in other words, it serves as a dark zoom lens. In addition, with a large number of glass lenses, the zoom lens is required to improve its cost and weight.\nIn conjunction with achieving further miniaturization of the liquid crystal projector, presently, achievement of its short-range projection, that is, its wider angle of view, which contributes especially to a home-theater-use projector as a great advantage is strongly requested.\nAlso, in order to achieve a higher brightness of a projected picture, the projection lens is required to be bright and have a large aperture ratio.\nIn general, when the projection lens is made so as to achieve a wider angle of view while maintaining a long back focus, its lens group closest to the magnification side has larger refractive power.\nAlso, in order to achieve a telecentric zoom lens on the reduction conjugate side, the overall lens group arranged from its aperture to the reduction conjugate side has larger positive refractive power, and the overall lens system has a retrofocus refractive-power arrangement. Hence, an asymmetry of the lens system increases, resulting in difficulty in compensating for especially distortion aberration, chromatic aberration of magnification, and the like.\nIn addition, a curvature of field increases without minimizing the Petzval sum in accordance with a wider angle of view, resulting in difficulty in compensating for this phenomenon."}
{"text": "SSL devices generally use semiconductor light emitting diodes (“LEDs”), organic light emitting diodes (“OLED”), and/or polymer light emitting diodes (“PLED”) as sources of illumination rather than electrical filaments, plasma, or gas. For example, FIG. 1 is a schematic cross-sectional diagram of a conventional indium-gallium nitride (InGaN) LED 10. As shown in FIG. 1, the LED 10 includes a substrate material 12 (e.g., silicon), N-type gallium nitride (GaN) 14, GaN/InGaN multiple quantum wells (“MQWs”) 16, and P-type GaN 18. The LED 10 also includes a first contact 20 on the P-type GaN 18 and a second contact 22 on the N-type GaN 14. During manufacturing, the N-type GaN 14, the GaN/InGaN MQWs 16, and the P-type GaN 18 are formed on the substrate material 12 via metal organic chemical vapor deposition (“MOCVD”), molecular beam epitaxy (“MBE”), liquid phase epitaxy (“LPE”), hydride vapor phase epitaxy (“HVPE”), and/or other epitaxial growth techniques, each of which is typically performed at elevated temperatures.\nOne operational difficulty of forming the LED 10 is that the N-type GaN 14, the GaN/InGaN MQWs 16, and the P-type GaN 18 may be delaminated from the substrate material 12 and/or otherwise damaged during high-temperature epitaxial growth and/or cool-down thereafter. Typically, the substrate material 12 includes silicon (Si), sapphire (Al2O3), silicon carbide (SiC), and/or other “non-native” materials because “native” materials (e.g., GaN or InGaN) with usable dimensions are difficult to produce. The non-native substrate materials have different coefficients of thermal expansion (“CTEs”) than the GaN/InGaN materials 14, 16, and 18. For example, the CTE of silicon is substantially less than that of GaN, and the CTE of sapphire is substantially greater than that of GaN. Such CTE differentials induce thermal stress as the wafer cools, which warp the substrate material 12 and/or cause crystal defects in epitaxial GaN/InGaN materials 14, 16, and 18. Additionally, the non-native substrate materials that facilitate particularly good epitaxial growth, such as Si(1,1,1) silicon wafer, can be expensive. Accordingly, several improvements in reliably and cost-effectively manufacturing SSL devices may be desirable."}
{"text": "1. Field of the Invention\nEmbodiments disclosed herein generally relate to the deposition of thin films using chemical vapor deposition (CVD) processing. More particularly, embodiments generally relate to methods for depositing barrier layers onto large area substrates.\n2. Description of the Related Art\nOrganic light emitting diodes (OLED) are used in the manufacture of television screens, computer monitors, mobile phones, other hand-held devices, or other devices for displaying information. A typical OLED may include layers of organic material situated between two electrodes that are all deposited on a substrate in a manner to form a matrix display panel having individually energizable pixels. The OLED is generally placed between two glass panels, and the edges of the glass panels are sealed to encapsulate the OLED therein.\nThe encapsulation is achieved by sealing the active materials in inert atmosphere using a glass lid secured by a bead of UV-cured epoxy resin. The rigid encapsulation is not applicable to flexible displays, where a durable and flexible encapsulation is necessary to protect the active OLED materials from water moisture and oxygen. One approach is to use a multilayer encapsulating structure as a barrier to water moisture and oxygen permeation. Inorganic layers can be incorporated into the multilayer encapsulating structure as the main barrier layer. Organic layers can also be incorporated for the purposes of stress relaxation and water/oxygen diffusion channels decoupling layer.\nSilicon nitride (SiN) is known as a good barrier material, thus it shows potential as an inorganic barrier layer in the multilayer encapsulation structure. However, SiN films deposited at low temperatures such as below 100 degrees Celsius have high stress, which can induce film peeling, also known as delamination, or mismatch issues in multi-film stack configurations. Due to the sensitivity of some of the layers of the OLED device, subsequently deposited layers over OLED materials will need to be deposited at lower temperatures, such as at temperatures less than 100° C.\nThus, there is a need for methods of depositing encapsulation/barrier films onto large area substrates with improved water-barrier performance to protect the devices underneath."}
{"text": "Field of the Invention\nThe present invention relates to an air duct and a cooling system for a vehicle. More particularly, the present invention relates to an air duct and a cooling system for a vehicle for improving cooling performance and aerodynamic performance.\nDescription of Related Art\nGenerally, an air duct is a passage through which air passes, and is a pipe for guiding air to parts requiring air from a part through which air can easily flow inside.\nFor example, there are an air duct that guides air to be sucked into an engine and an air duct for cooling a brake system.\nRecently, techniques for preserving temperature of an engine and minimizing fuel consumption at the time of initial starting of the engine have been actively developed.\nHowever, cooling of an engine may not be performed well, fuel consumption may be deteriorated during high speed driving, and heat damage may be generated to components which are disposed to a periphery of an exhaust passage at a high temperature if only the method for preserving temperature of an engine is performed. In addition, the performance of cooling an engine and fuel consumption may be deteriorated, and the heat damage may be become serious if air flowing through a radiator grille is interfered with by ancillary machinery disposed to the periphery of an engine.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an image sensor, and more particularly to an image sensor that performs a Time Delay Integration (TDI) scanning at high speed.\n2. Description of Related Art\nA charge coupled device (CCD) line sensor is a one-dimensional sensor in which photo-sensors receiving image light are arranged in line. To image a two-dimensional image, the CCD line sensor or the object is moved so that the object can be imaged on a line-by-line basis. A scanner having the CCD line sensor is used in a copy machine and a product inspection machine, for example.\nWhen the CCD line sensor captures the object that is moving at high speed, or when the CCD line sensor is moved at high speed to capture the object, the signal charge must be repeatedly accumulated and transferred at high speed on a line-by-line basis. Consequently, the signal charge can be accumulated for only a short period of time per line. For this reason, there is a problem in that a sufficient quantity of light is not obtained for imaging.\nTo solve this problem, a TDI sensor is used to scan the object at high speed. The TDI sensor is constructed in such a way that a plurality of CCD line sensors are arranged in a scan direction. The TDI sensor transfers signal charge accumulated in the CCD of each line to the CCD of the next line in synchronism with the movement of the image. Consequently, the signal charge is successively accumulated in the plurality of CCD line sensors, and the high-speed scanning can obtain a sufficient quantity of light for imaging.\nThe scan speed of the TDI sensor is limited. That is because there is a limit to a speed at which the signal charge accumulated in the CCD is shifted to the outside. Specifically, the signal charge accumulated in the CCD line sensors is finally discharged to the outside as serial analog signals, and an A/D converter converts the analog signals into digital signals. Owing to the limited processing capacity of the A/D converter, the scan speed is limited.\nTo raise the scan speed, one CCD line sensor is divided into a plurality of patches, and the signal charge accumulated in each patch is discharged as serial analog signals, which are converted into digital signals by each of a plurality of A/D converters provided for each patch. In this case, however, there is a limit to the width that can be divided into patches, and therefore, the scan speed cannot be raised sufficiently.\nThe present invention has been developed in view of the above-described circumstances, and has as its object the provision of an image sensor that realizes a high-speed scanning.\nTo achieve the above-mentioned object, the present invention is directed to an image sensor comprising: a light receiving part having photo-sensors, the photo-sensors receiving object light and converting the object light into signal charge, the photo-sensors being arranged in a plurality of lines in a line direction and being arranged in a plurality of lines in a scan direction perpendicular to the line direction; a plurality of scan direction signal charge transfer parts for transferring the signal charge, read from the photo-sensors, in the scan direction in accordance with a scan speed and accumulating the signal charge representing the same image read from the photo-sensors in the same line, each of the scan direction signal charge transfer parts being provided for each line of the photo-sensors in the scan direction; a plurality of line direction signal charge transfer parts for transferring the signal charge, transferred through the scan direction signal charge transfer parts, in the line direction, each of the line direction signal charge transfer parts connecting with every a predetermined number the scan direction signal charge transfer parts; a plurality of digitizing parts for converting a plurality of signal charge lines, transferred through the line direction signal charge transfer parts, to digital signals; and an image signal producing part for producing image signals according to the digital signals output from the digitizing parts.\nAccording to the present invention, the signal charges are read from the photo-sensors and are transferred in the scan direction in accordance with the scan speed, and the signal charges read from the photo-sensors on the same line, which represent the same image, are accumulated. The signal charges of every the predetermined number lines in the line direction are sequenced in a signal charge line in the line direction. The signal charge lines are transferred in the line direction, and are separately converted to the digital signals, from which the image signals are generated.\nThus, the signal charge line representing the image of one line can be converted into the digital signal by a plurality of processing circuits. For this reason, one processing circuit processes only a small amount of data, and the object can be scanned at high speed.\nThe photo-sensors are arranged at intervals of more than predetermined distance at the light receiving part, and the light receiving surface of each photo-sensor can be enlarged by shifting the peripheral circuits. As a result, the light receiving surface can receive a large quantity of light."}
{"text": "The invention relates to handles for carrying batteries, and more particularly to a rope-type battery carrying handle that has an end of the rope removably attached to an end of a grip.\nStarting, lighting, and ignition (SLI) batteries are typically used in automotive, recreational, and other applications, are heavy, cumbersome, and usually require two hands, or often two people, for carrying. The desirability of providing such batteries with attachable/detachable handles for facilitating carrying, placement, and retrieval of such batteries has long been known. Such handles are a particular convenience in batteries designed for use in boats or in uninterrupted power supply (UPS) applications which must be frequently moved for storage, service, or recharging.\nBail-type handles, which are known in the art, typically comprise a U-shaped or C-shaped member attached to opposing sides of a battery casing, either on its container or cover. With such handles, the battery may be carried in much the same fashion as a picnic basket or bail.\nSubstantially rigid bail-type handles are known in the art. A variety of such handle designs have been proposed for carrying batteries. Detachable, substantially rigid bail handles are disclosed, for example, in U.S. Pat. No. 3,093,515 to Rector, U.S. Pat. No. 3,956,022 to Fox, U.S. Pat. No.4,029,248 to Lee, U.S. Patent No. 4,673,625 to McCartney et al., U.S. Patent No. 5,232,796 to Baumgartner, U.S. Pat. No. 5,242,769 to Cole et al., and U.S. Pat. Des. No. 292,696 to Sahli.\nRope-type handles are likewise known in the art. Rope-type handles typically have one or more injection molded plastic part coupled by flexible rope sections and, accordingly, are physically highly flexible. The rope sections are generally a braided synthetic material such as polypropylene.\nAccording to one type of rope handle design, the ends of the rope handle are manually fed into two holes and coupled to the battery container. In the battery disclosed in U.S. Pat. No. 3,092,520 to Buskirk et al., the rope handle is coupled to the battery container by cementing the ends of the rope in recesses in projections on the sides of the battery container. Alternately, the ends of the rope handle may include an enlarged molded plastic portion and may be pressed into slots underneath the handle bracket area as shown, for example, in U.S. Pat. No. 3,797,876 to Gummelt and U.S. Pat. No. 4,013,819 to Grabb. According to other designs, the ends of the rope may be enlarged as shown for example in British Patent 869,329, or the ends coupled or welded together as shown for example in British Patent 869,329 and British Patent 1,453,977.\nIn another type of rope handle design, looped rope portions extend from the ends of a molded plastic grip portion as shown, for example, in U.S. Pat. No. 971,876 to Apple, U.S. Pat. No. 4,791,702 to McVey, and U.S. Pat. No. 5,242,769 to Cole et al. The looped rope portions are then coupled to the battery container via dedicated protrusions extending from the walls of the battery by looping the rope around the protrusion and then securing it into a recess or the like.\nAnother such rope handle design is disclosed in U.S. Pat. No. 5,144,719 to Arthur. The Arthur patent discloses a xe2x80x9cU-shapedxe2x80x9d handle having one end of the rope embedded in one depending leg of the handle. The opposite end of the rope includes an enlarged head, which may be fed through lugs on the battery. The enlarged head and the adjacent length of rope are then laid into a tri-part vertical slot on the other depending leg of the handle, the head being disposed in the upper portion of the slot, the adjacent rope extending through the lower two portions of the slot. Significantly, however, the head and adjacent rope section are not secured to the handle. As may be seen in the illustrations of the reference, there is sufficient clearance between the head and the slot, as well as the adjacent rope section and the slot such that the head and rope section may become easily dislodged from the handle leg unless a constant vertical force is maintained on the handle. Accordingly, the Arthur handle does not provide an attachment mechanism which is reliable. Moreover, the intricate coupling design requires the user to have a high level of manual dexterity and a working knowledge of the defailed structure of the complex attachment.\nInstallation of these rope handle designs may be labor intensive. Properly securing the ends of the rope to the battery container or securing the loop ends around a protrusion and into a recess can be quite time consuming and may require manual dexterity. These difficulties in installing the battery handles can lead to improper installation, which can result in an unreliable battery handle.\nAdditionally, these designs generally require specialized handle brackets to be molded into specific containers. Complicated grip and/or rope end configurations may also be required. These requirements can result in increased costs in the form of mold and tooling costs, as well as increased labor and downtime costs during changeover. Further, storage and floor space costs increase because the battery manufacturer must maintain larger inventories.\nIt is a primary object of the invention to provide a rope handle that may be reliably and easily assembled onto a battery container and which remains securely coupled to the battery until purposely removed by the user.\nA related object of the invention is to provide a rope handle arrangement that has a relatively simple design, and does not require high manual dexterity to assembly for a secure, reliable handle.\nIt is a further object of the invention to provide a rope handle that may be utilized with a battery that produces an acceptable appearance.\nIt is another object of the invention to provide a rope handle that contributes to the production of an economical battery. A related object of the invention is to provide a rope handle design that minimizes manufacturing and inventory costs.\nThese and other objects and advantages of the invention will be apparent to those skilled in the art upon reading the following summary and detailed description and upon reference to the drawings.\nIn accomplishing these and other objects of the invention, there is provided a battery that includes rope handles each of which engages a handle bracket on an end wall of a conventional battery container. Each rope handle includes a grip with a retaining recess at one end of the grip and a rope secured to the other end of the grip by molding or the like. The rope has an enlarged distal end or a cylindrical plug molded for engaging the retaining recess of the grip. The retaining recess includes a generally keyhole-shaped slot which extends through the grip from a first surface to a second surface and which has a hole portion and a channel portion projecting radially from the hole portion and terminating at an end. The retaining recess also includes a counterbore located on the first surface of the grip and encompassing the hole portion of the slot. To secure the rope to the grip, the rope is slid through the slot and the plug is subsequently drawn towards the grip and is retained within the counterbore, thus securing the handle to the battery container. In other words, the retaining recess includes a counterbore with a subjacent retaining surface for receiving and supporting the plug, and radially extending slot. The rope is laterally advanced through the slot to move the plug into position above the counterbore. The plug is then pushed down into position in the counterbore and/or a downward force is exerted on the rope to position the plug and secure the rope handle."}
{"text": "Immobilized Reagents For Removal Of Circulating Immune Reactants In Vivo\nThe role of immune reactants in many experimental and human diseases is now well established, Cochrane C. C., Koffler D.: Immune complex disease in experimental animals and man. Advances Immunol 16:185-233, 1973; Wilson C. B., Border W. A., Lenham D. H.: Renal diseases in Basic and Clinical Immunology. (Fudenberg, HH, ed.). Lange Publications, Los Altos, Calif., 1976, pg. 562. Therapy for many of these immunologically mediated diseases has consisted largely of the use of pharmacological agents that widely and non-specifically suppress host immunity leading to numerous undesirable effects. With increasing awareness of the etiopathogenic factors in many autoimmune diseases many sensitive radioimmunological techniques to measure them in serum. Various immunoadsorbents designed to extract pathogenic immune reactants from the circulation have been developed. For example, immunoadsorbents consisting of immobilized antigens, antibodies and enzymes have been developed. When placed in an extracorporeal circuit, these immunoadsorbents have shown a capacity to extract or hydrolyze immune reactants in the circulation without demonstrable release of immobilized substances or significant immediate or long-range toxicity to the host.\nGraf et al, Graff M. W., Uhr J. W.: Regulation of antibody formation by serum antibody. I. Removal of specific antibody by means of immunoadsorption. J Exp Med 50:130-1175, 1969, were the first to show that immunoadsorbents could be employed to selectively remove antibodies from actively and passively immunized rabbits. In later studies, Schenkein et al., Schenkein I., Brystryn J. C., Uhr, J. W.: Specific removal of in vivo antibody by extracorporeal circulation over an immunoadsorbent in gel. J Clin Invest 50:1864-1870, 1971, developed an extracorporeal immunoadsorbent system in which bovine serum albumin (BSA) was immobilized in agarose and proved capable of selectively removing BSA antibodies from the circulation. A similar immunoadsorbent in which ssDNA antigen was immobilized and proved capable of extracting ssDNA antibody in both passively and actively immunized rabbits, Terman D. S., Stewart I., Robinette J., Carr R., Harbeck R.: Specific removal of DNA antibodies in vivo with an extracorporeal immunoadsorbent. Clinical and Experimental Immunology 24:231-238, 1976. Because of the fragility of the supporting matrix and the possibility of leaching of immobilized substances, new and more stable extracorporeal immunoadsorbents were subsequently developed.\nVarious antigens such as bovine serum albumin (BSA), deoxyribonucleic acid (DNA), glomerular basement membrane (GBM) extract have been immobilized on several solid supports and these have been employed as extracorporeal immunoadsorbents, Terman D. S., Durante D., Buffaloe G., McIntosh R.: Attenuation of canine nephrotoxic glomerulonephritis with an extracorporeal immunoadsorbent. Scandinavian Journal of Immunology 6, 1977: Terman D. S., Petty D., Ogden D., Harbeck R., Buffaloe G., Carr R.: Specific removal of DNA antibodies in vivo by extracorporeal circulation over DNA immobilized in collodion-charcoal. Clinical Immunology and Immunopathology 8, 1977; Terman D. S., Tavel T., Petty D., Racic M. R., Buffaloe G.: Specific removal of antibody by extracorporeal circulation over antigen immobilized in collodion-charcoal. Clinical and Experimental Immunology 28, 1977; Terman D. S., Tavel T., Petty D., Tavel A., Harbeck R., Buffaloe G., Carr R.: Specific removal of bovine serum albumin (BSA) antibodies by extracorporeal circulation over BSA immobilized in nylon microcapsules. Terman et al, Journal of Immunology 116:1337, 1976; Terman D. S., Stewart I., Robinett J., Carr R., Harbeck, R.: Specific removal of DNA antibodies in vivo with an extracorporeal immunoadsorbent. Terman et al, Clin Exp Immunol 24:231, 1976.\nCirculating immune complexes have now been implicated in the pathogenesis of numerous diseases. Their presence in the circulation often correlates with disease activity and they may be found deposited in tissues, Zubler R. H., Lambert P. H.: Detection of immune complexes in human diseases. Prog Allergy 24:1, 1978. In addition to studies described above for hydrolysis of nDNA:antiDNA complexes, preliminary work has shown that Clq, the first component of complement, may be immobilized in collodion membranes and will bind to immune complexes circulated over them. Terman et al FEBS, Letters, 68, 89, 1976."}
{"text": "The thermoregulatory system of homeotherms has an inherent ability to hold core body temperature within a small variation of a set point. Excursions above and below the set point can cause compromised body function, injury, and even death may occur.\nOperation of the thermoregulatory system is based on a complex, nonlinear network of feedback control signals and responses to adjust the thermal resistance between the body core and the environment and to modulate the rate and distribution of internal energy generation. The operation of the thermoregulatory system is remarkably efficient over a broad spectrum of physiological states and environmental conditions.\nIn certain circumstances, however, the thermoregulatory system is unable to maintain the core temperature within the set operational range, or there may be therapeutic or prophylactic reasons to override the system to maintain a reduced core body temperature."}
{"text": "Electrical cables for transmission of electrical signals are known. One common type of electrical cable is a coaxial cable. Coaxial cables generally include an electrically conductive wire surrounded by an insulator. The wire and insulator are typically surrounded by a shield, and the wire, insulator, and shield are surrounded by a jacket. Another common type of electrical cable is a shielded electrical cable comprising one or more insulated signal conductors surrounded by a shielding layer formed, for example, by a metal foil. To facilitate electrical connection of the shielding layer, a further un-insulated conductor is sometimes provided between the shielding layer and the insulation of the signal conductor or conductors."}
{"text": "Modern organizations, particularly technology driven organizations, have a constant need to innovate. Unfortunately, innovation is often hindered by an existing lack of technology to properly explore new technologies. In a sense, future technologies are “stacked” upon previous technologies. That is, without a proper technological foundation, there is no support for technological innovation. It is difficult to skip generations of technology.\nOne technology often needed to support innovation is data processing power. Modem organizations require processing power for many different purposes. Technology companies, for instance, rely on processing power for research and development (“R&D”) efforts. Pharmaceutical research companies spend large sums of money on the latest data processing equipment to discover and test new drug compounds. Financial institutions need processing power for stock trend modeling and other computation-intensive activities. Defense contractors need processing power to test new missile designs. While these various organizations operate in different industries, each would benefit by owning or having access to computer systems that offer more data processing power.\nAt any given time, there are millions of unused data processing devices throughout the world. Workstations, personal computers, laptop computers, personal digital assistants (“PDAs”) are often not in use. Many computers sit in a room overnight, during meals, and sometimes throughout the workday. Further, even when people are using these data processing devices, the devices rarely operate at full capacity. The processors in these data processing devices are often available to provide data processing for others.\nData processing devices are possessed by various entities including corporations, academic institutions, and private individuals. These entities, however, are generally only concerned with solving their particular problems. The entities described above have no incentive to share the processing power of their respective devices with others. Thus, most never consider the greater need for data processing power."}
{"text": "The present invention relates generally to a device for detecting whether any one of a number of fuses in an electrical system has blown out.\nA machine, such as an automobile, which incorporates a large number of electrical devices, such as ignition loads, radio sets, various different lamps, etc., requires a relatively large number of fuses for protecting the loads from excessive current, and for guarding against fire and so on caused by short circuits. Various devices have been proposed for detecting whether any of thses fuses have been blown out. For example, a fuse blowing detector such as shown in FIG. 1 of the accompanying drawings has been suggested. In this device, an automobile has an electric system which includes loads L1, L2, L3 fed from a battery 1 through an accessory switch SW1, loads L4 and L5 fed through an ignition switch SW2, and loads L6 and L7 fed directly from the battery. These loads L1-L7 are fused by fuses F1-F7 respectively. Several switches are shown as controlling supply of power to some of the loads, such as S1, S3, S4, and S6. The fuse blowing detector shown in FIG. 1 includes wires to each of the junctions between a fuse and its load, and a slidable changeover switch 2 which has a plurality of fixed contact points 2a connected to the other ends of these wires. The slidable contact 2b is capable of contacting each of these fixed contacts, by the movement of a slide lever 3 and an indicator lamp 4 having one terminal connected to the slidable contact 2b and an other terminal to the earth of the vehicle.\nIn operation, the ignition switch SW2 and the accessory switch SW1 are closed, and the slide lever 3 is moved to and fro. If all the fuses are intact, the lamp 4 is lighted each time slidable contact 2b and a fixed engage 2a contact. However, if any one of the fuses is blown, contact between slide lever 3 and the corresponding fixed contact, does not cause lamp 4 to light, thereby signalling blowing of the fuse.\nThis detector, however, has certain disadvantages because the use of a mechanical part, such as the slide lever 3, makes the detector prone to malfunctions. Another disadvantage is that the number of movements of the slide lever 3 for each test sequence is equal to the number of fuses. Further, use of only one lamp has made it difficult to know which fuse is defective."}
{"text": "Two critical parameters effect the quality and consistency of a weld between two or more work pieces, weld gun pressure and the current density in the region to be welded. Weld gun pressure is readily regulated by pneumatic, hydraulic, or other mechanical means. Current density regulation requires an electronic solution. Many methods have been utilized to regulate and maintain the current density constant within the contact area between weld gun contact tips and the material to be welded. As the contact tips deteriorate, the contact area increases, resulting in a decrease in the current density at the weld nugget. This results in a decreased heat input and can result in weld defects. Compensation for this decrease in current density over the life of the tips can be accomplished through several different methods to increase or boost the current. Less heat and thus less current, is required during the first or early stage of the contact tips' life. Once the contact tips have settled in, during a second stage, a gradual increase in heat is required. During the last stage, as the contact tips start to deform, even more heat is required. These three stages form a user profile for the current or heat boost.\nEarlier weld controllers modified the firing angle of SCR switches to regulate the conduction angle of the SCRs to a particular percentage of full or 180 degree conduction. This mode of operation, known as a voltage compensation method, does not regulate current directly. A second method measures the available heat as a function of the overall system power factor, and provides a user programmed percentage of that available heat following a profile based on the above mentioned three stages of the contact tip lifetime. This method will or could provide for line voltage variations. The user adjusts either the percent conduction angle or the percent heat in either of these methods to achieve a desired weld current as measured by external means.\nSome weld controllers provide a third method which uses a constant current control which will adjust the firing angle of the SCRs to maintain a predetermined current flow to the contact tips based on the user profile. The use of stepper programs implements this method by increasing the current in equal increments according to the user profile. Some prior art weld controllers employ a manual stepper to adjust for the current boost, which typically is increased as a series of scheduled linear steps as specified by a weld engineer to obtain metallurgically sound welds during the life of the weld contact tips. For example, the first stage may be programmed to reach a 5 percent current boost in one percent increments after 200 weld cycles, the second stage may be programmed to reach a 10 percent current boost after 2000 weld cycles, and the last stage may be programmed for a 15 percent current boost after 8000 weld cycles. Commonly assigned U.S. Pat. No. 5,083,003 discloses an adaptive stepper which increases the heat boost and thus the current density as a function of not only the number of weld cycles but also as a function of expulsions. Expulsions, also known as spitting, generally indicate that too much heat is being applied during the weld cycle. Molten material is blown away from the weld zone during expulsion, resulting in a significant drop in resistance at the primary of the weld transformer supplying the contact tips.\nIn all of these cases, it becomes difficult to detect process variations in the welding cycle that could indicate other fault conditions. These variations occur as a result of either short term or long term impedance changes in the welder system. Short term changes are caused by variations in the workpiece and contact tip interface such as oxidation of the surface of the workpieces or poor part fit Long term effects are associated with tool or weld cable degradation or poor connections. With the first and second methods, this increased impedance will result in lower weld current. In the third method, the weld controller will attempt to directly compensate for the increase and regulate the current to maintain it at the constant, preset level, providing more and more power to the weld system, with possible expulsions occurring. The user would have no indication of a problem until it may be too late, as the system would have to detect an inability to regulate current to reach a current limit before error messages are generated. At that time, it would only indicate a failure of the control to regulate the weld current, which could be from many sources. There would be no indication of a change in system impedances. The current controlled system will continue to make metallurgically sound welds until it fails catastrophically from lack of maintenance. A method to indicate changes in system impedances is therefore desirable.\nQuality and strength of a spot weld can be correlated with a change in resistance as measured through the weld as the weld progresses during the fusion process. This will effect a change in the power factor of the circuit which will be reflected from the secondary circuit of the welding transformer coupling power to the contact tips to the primary circuits and will result in a significant drop in resistance at the primary circuit. The timing changes resulting from this change can be sensed by the microprocessor based weld controller system. The amount of increase or decrease in the current conduction angle can be determined from these changes and can become a basis for controlling the welding heat applied to the workpiece being welded. U.S. Pat. No. 4,399,511 describes one such type of control. The controller measures the weld current conduction time from the point of initiation to the point of extinction. The conduction time, when added to the original time delay of the initiation signal, is related to the power factor. Using a numerical representation a characteristic resistance curve that relates the ohmic resistance of the work piece as a function of time and weld cycles, the measured and calculated power factor changes can be compared with this curve to determine if more or less energy is required to be supplied to the welding tips. This requires the use of look-up tables or complicated calculations that are non-trivial to calculate the power factor for each weld cycle."}
{"text": "Semiconductor devices are used in a variety of electronic applications, such as personal computers, cell phones, digital cameras, and other electronic equipment, as examples. Semiconductor devices are typically fabricated by sequentially depositing insulating or dielectric layers, conductive layers, and semiconductive layers of material over a semiconductor substrate, and patterning the various material layers using lithography to form circuit components and elements thereon.\nThe semiconductor industry continues to improve the integration density of various electronic components (e.g., transistors, diodes, resistors, capacitors, etc.) by continual reductions in minimum feature size, which allow more components to be integrated into a given area. These smaller electronic components also require smaller packages that utilize less area than packages of the past, in some applications.\nOne smaller type of packaging for semiconductors is a flip chip chip-scale package (FcCSP), in which a semiconductor die is placed upside-down on a substrate and bonded to the substrate using bumps. The substrate has wiring routed to connect the bumps on the die to contact pads on the substrate that have a larger footprint. An array of solder balls is formed on the opposite side of the substrate and is used to electrically connect the packaged die to an end application.\nHowever, some FcCSP packages tend to exhibit bending, where warping of the substrate occurs during processing, such as during temperature stress. The bending can cause reliability issues, such as bond breakage of the bumps, delamination of an underfill, and delamination of a passivation layer on the die."}
{"text": "In a tube and shell heat exchanger, such as a vertical tube chemical reactor, there are various situations in which it is desirable to temporarily plug or indicate the condition of one or more tubes. A temporary plug may be used to protect a tube or catalyst inside the tube, or to identify the condition of the tube, or to aid in keeping track of the work progress on tubes while the reactor is out of service for maintenance. In the prior art, this is usually done by a colored plastic cap or plug inside the top end of each tube, with a particular cap or plug color intended to identify a particular tube condition (such as a tube to be unloaded, cleaned, loaded, pressure-drop tested, tube failed due to high pressure drop, tube failed due to a low pressure drop, tube passed the pressure drop test, tubes that have or will have thermocouples or pressure sensors, or even tubes from which catalyst samples may be removed for laboratory analysis after the catalyst has been used).\nThis procedure requires large numbers of caps or plugs of different colors to be used as the tubes in the reactor are worked on to indicate the status of particular tubes and to serve as a visual indication of the work flow progress on all of the tubes in the reactor while the tubes are brought to within the desired specifications for cleaning, empty, full, outage, and pressure drop. It is not unusual to run out of different colors to designate the numerous tube conditions, resulting in improvisations by the catalyst changeover crew in order to identify the condition of the tubes. For instance, a cap or plug may be removed from a tube and if pliable enough then twisted and reinserted sideways into the top edge of the tube to designate the new condition of the tube. Sometimes two caps or plugs are stacked on top of each other and inserted into a tube. In both of these cases, the caps or plugs then project awkwardly from the tube and are more prone to being accidentally kicked and dislodged from the tube, so they no longer perform the desired function of visually identifying the tube condition. In the case of the twisted and reinserted cap or plug, it can become lodged down inside the tube if walked upon, such that it may not be easily seen and may even unknowingly be left behind.\nIn order to prevent the caps or plugs from being accidentally dislodged from, or pushed down into the tubes, they usually are designed to fit snugly inside the tubes. Some caps or plugs are simple pipe thread protectors that are applied to protect the threads on the ends of threaded pipes that use a tapered thread like those available from Caplugs from Buffalo, N.Y. (www.caplugs.com). Since these devices are tapered, they can be pushed very securely into a tube. Upper tubesheets with caps or plugs in them are continuously walked upon by the catalyst change crew, and heavy machines and supplies are temporarily located and moved across this upper tube sheet, forcing the caps or plugs securely into the tubes.\nCap or plug size availability is often somewhat limited, as these caps or plugs are sized for commercially sized pipe threads, although custom sizes could be used. Reactor tube inner diameters are seldom standard pipe sizes, and instead are often custom sized to optimize heat transfer. In the case of tight fitting tapered caps or plugs, this results in two problems. First, the tight wedging action of the tapered cap or plug within or upon the tube makes it very difficult to remove the cap or plug from the tube, so it usually is pried but with a screwdriver, which may damage the cap or plug and may even score the inside wall of the tube. Second, after a cap or plug which has been removed and reinserted a number of times may become; cut or otherwise damaged by the blade of a screw driver, or by being removed by narrow end (needle nose) pliers, it may no longer fit snugly inside the tube, making it more prone to being accidentally kicked out or otherwise, dislodged from its tube, or it may be damaged so that it is no longer serviceable.\nTape, such as duct tape, can be attached to the caps or plugs in certain tubes in order to identify them (this tape may be accidentally scuffed off or otherwise removed from its corresponding cap or plug). Sometimes caps or plugs are actually taped to the tubesheet (resulting in art undesirable sticky residue on the tubesheet when the tape is removed).\nMost tube caps or plugs, such as those intended to protect the threads on the ends of pipes, include an open cavity or depression where dirt, catalyst and foreign material can accumulate. It can be difficult and time consuming to remove the dirt, spilled catalyst and foreign material from these cavities, especially when many thousands of such caps or plugs are used.\nSome tube plugs have a flat top (do not include an open cavity) but they incorporate a tapered body that fits snugly into the tubes and require a screw driver or fork device to remove them, making them time consuming and sometimes difficult to install and remove. Other plugs have a flat top but have a loose fitting body so they effectively float around in the tube. These plugs are relatively easy to install and remove, but they do not seal the tube from air flow, and they are sometimes inadvertently removed while walking across them or by moving equipment across them. Since reactors are often made of carbon steel, the reactor is often kept warm to prevent condensation that could lead to iron oxide formation. Also, it is desirable to keep most catalyst pellets that have been loaded dry and away from moisture. The warm vertical tubes in the reactor vessel may induce natural convection air flow through the tubes, and this may be undesirable as the air being supplied to the bottom of the tube could be moisture-laden ambient air. Another reason for minimizing or eliminating this natural convection phenomenon is that some catalyst can off-gas certain chemicals which can be annoying and even hazardous for the catalyst handling crew. Also, convection air flow can raise the ambient temperature in the dome area above the tubesheet, making it uncomfortable for the catalyst handling personnel. Loose fitting plastic caps with or without recesses or plugs with flat tops can slow down the convection air flow, but they do not stop it."}
{"text": "In the field of dental implants, patient comfort and the efficient use of a dentist's time are paramount. Likewise, precision alignment of the prosthetic components is essential. The need to match upper and lower teeth to within a few microns and provide accurate mating of the prosthesis with existing teeth requires accurate replication of oral structures when making dental impressions. It is possible with the introduction of an intra-oral scanning apparatus to scan the oral environment and precisely display a working virtual model on a computer screen and to generate a physical model in accurate detail. Alternately, the more direct method of creating a physical stone model is still preferred by the majority of dentists.\nTo create accurate stone models for fitting the final prosthesis, a matched upper and lower impression can be obtained simultaneously with the jaw in the closed position and the teeth in the interdigitated position (centric occlusion). Currently, the tall impression transfer posts used to register the implants to the upper and lower jaws prevent the full closure of the mouth while making the simultaneous upper and lower scans and impressions. The present invention remedies this oversight for both virtual scanning and physical modeling.\nExisting practice has been to perform the following procedures. After dental implants have healed into the underlying bone structures of the mandible or maxilla and the soft tissue has healed around a protective healing cap or healing screw, a full set of upper and lower impressions of the mouth are made using individual full or partial arch upper and lower trays. A separate bite registration elastomeric impression is also made in centric closure. Positive casts of the upper and lower impressions are poured, mounted on a dental articulator and aligned in centric occlusion by means of the separate elastomeric bite registration.\nIn more detail, the prior art performs the following steps to make an accurate impression, the healing caps are removed from one or more dental implant fixtures and standard length impression transfer posts are accurately placed with retaining screws on each implant fixture. An impression tray filled with a self-hardening elastomeric impression material is pressed over the region of the dental arch containing the impression transfer posts.\nAfter a few minutes, the elastomeric impression material has set. The tray is gently removed from the mouth leaving the tall standard impression posts still attached to the implant fixtures. The screw retaining the standard impression transfer post is removed and the tall impression post is attached to an implant analog. The impression post and attached implant fixture analog is reinserted in the elastomeric compound. Taking care to seat the impression post accurately. The pratitioner or a dental lab casts a stone model embedding the implant analog. The implant fixture analog is thus fixed accurately within the stone model. The stone model serves as the platform to craft the prosthesis. The healing cap is replaced on the dental implant fixture in the mouth. Another impression of the mating jaw is taken by the same method. A third elastomeric cast is made of the teeth in centric closure to insure later alignment.\nThe stone models of the upper and lower mouth structure with dental implant analogs exactly aligned and retained are molded in hard plaster stone from the separate impressions. These models are separately placed upon a dental articulator to mimic the actual jaw motions. The separate upper and lower stone casts are aligned in centric closure with the elastomeric bite cast. The final prosthesis is built and tried in for a non-interfering, good fit. This prosthesis relies upon proper replacement of the tall standard impression post to insure the properly aligned position.\nWith the optical scanning apparatus, a virtual three dimensional image of the whole oral environment can be created in a minute or two. The virtual image file can be transmitted electronically to a lab and a physical model of the upper and lower jaw can be printed in rigid polymer with a three dimensional printer.\nApplicants offer an impression post serving the needs of both the physical casting techniques and the virtual scanning methods.\nApplicants, in order quickly to make an accurate, simultaneous impression of the upper and lower teeth in the correct alignment use a dual arch impression tray such as the Triple Tray™. This tray consists of a molded plastic or metal assembly with a handle connected to a set of confining dams and a thin open screen mesh. The mesh is oriented horizontally and is to be placed between the mating occlusal surfaces of the teeth while the jaw is in the closed or centric position. The mesh is thin enough to allow complete centric closure. The mesh is flexible and porous to trap and retain an elastomeric impression compound.\nThe buccal and lingual dams are molded to the mesh. A paste of quick-setting elastomer is placed on both sides of the mesh within the confines of the dams. The mouth is closed with the upper and lower teeth in the closed or centric position while embedded within the curing elastomer. In this manner, a matching set of aligned upper and lower impressions along with the proper bite registration are made simultaneously.\nThe elastomeric impression materials, such as polyvinylsiloxane or polyether, are dimensionally stable, but need adequate thickness and surface area in contact with the impression transfer post to ensure accurate positioning and replication of the implant within the models mounted upon an articulator. Currently, long tapered impression transfer posts are used, which have adequate surface area to accurately register the elastomeric impression to the dental implant analog, but interfere with the use of a dual arch impression tray by preventing closure.\nIn prior art Neal B. Gittleman U.S. Pat. No. 6,213,773 Reduced Height Dental Impression Post, and Neal B. Gittleman U.S. Pat. No. 7,632,096, with all reference made therein, a low profile wing impression post is taught. The impression post is of reduced height to allow complete closure to obtain a triple impression of the upper and lower dentition as well as the dental occlusion. This impression post is removably attached to an installed implant fixture and remains embedded within the triple impression compound, when the impression is removed from the mouth. The winged coronal top of the impression transfer post has sufficient volume and surface area to remain accurately fixed and stable within the impression compound. An implant fixture analog is mated with the impression transfer post and cast in the stone model. The prosthesis is accurately built upon implant fixture analog, with a matching manufactured abutment."}
{"text": "The present invention relates to a thin film type piezoelectric element which transduces electric energy to mechanical energy and vice versa. The piezoelectric element is used as a pressure sensor, temperature sensor, actuator for ink jet recording head or the like. The present invention also relates to such an ink jet recording head. More particularly, the present invention relates to a process for the preparation of the thin piezoelectric film element.\nIn conventional ink jet recording heads, the vibrator which acts as a driving source for injecting an ink is composed of a thin piezoelectric film element. This thin piezoelectric film element normally comprises a thin piezoelectric film made of a polycrystalline substance and an upper electrode and a lower electrode arranged with the thin piezoelectric film interposed therebetween.\nThis thin piezoelectric film is normally made of a binary system having lead zircotitanate (hereinafter abbreviated as xe2x80x9cPZTxe2x80x9d) as a main component or a tertiary system comprising the binary system having a third component incorporated therein. The thin piezoelectric film having such a composition may be formed, e.g., by sputtering method, sol-gel method, laser abrasion method, CVD method or the like.\nA ferroelectric material comprising a binary PZT is disclosed in xe2x80x9cAllied Physics Lettersxe2x80x9d, 1991, Vol. 58, No. 11, pages 1161-1163.\nFurther, JP-A-6-40035 (The term xe2x80x9cJP-Axe2x80x9d as used herein means an xe2x80x9cunexamined published Japanese patent applicationxe2x80x9d) and xe2x80x9cJournal of The American Ceramic Societyxe2x80x9d, 1973, Vol. 56, No. 2, pages 91-96 disclose a piezoelectric material comprising a binary PZT.\nIn the case where a thin piezoelectric film element is applied to an ink jet recording head, it is preferred that a thin piezoelectric film (PZT film) having a thickness of from about 0.5 xcexcm to 25 xcexcm be used. This thin piezoelectric film must have a high piezoelectric strain constant.\nIn general, it is reportedly necessary that the PZT film be subjected to heat treatment at a temperature of 700xc2x0 C. or higher to allow the crystal grains in the thin piezoelectric film to grow in order to obtain a thin piezoelectric film having a high piezoelectric strain constant. As the material constituting the lower electrode in the thin piezoelectric film element there may be used an electrically-conductive material such as platinum, titanium, gold and nickel.\nJP-A-6-116095 describes crystal grains constituting a piezoelectric material. This patent discloses a process for the formation of a thin ferroelectric film which comprises applying a precursor solution of lead zircotitanate or lanthanum-containing lead zircotitanate to a platinum substrate which is oriented in (111) plane, characterized in that the-application of the precursor solution is followed by heat treatment at a temperature of from 150xc2x0 C. to 550xc2x0 C., where a desired crystalline orientation is attained, further followed by calcining at a temperature of from 550xc2x0 C. to 800xc2x0 C. for crystallization, whereby a specific crystal plane of the thin film is preferentially oriented along the surface of the substrate according to the heat treatment temperature.\nAs a prior art technique concerning the present invention there is proposed a process for the preparation of a bulk piezoelectric ceramic as disclosed in JP-A-3-232755. As disclosed in this reference, it is said that a piezoelectric ceramic having a higher density exhibits better piezoelectric characteristics.\nFurther, JP-A-50-145899 discloses an example of the application of a bulk piezoelectric ceramic in the generation of a high voltage as in gas apparatus. This patent describes that a piezoelectric ceramic having pores having a diameter of from 4 to 10 xcexcm uniformly dispersed therein and having a specific gravity of from 90% to 93% based on the true specific gravity exhibits a percent discharge rate of 100%.\nA conventional ink jet recording head comprising a thin piezoelectric film element is proposed in, e.g., U.S. Pat. No. 5,265,315.\nIn the case where a thin piezoelectric film (PZT film) having a thickness of not less than 1 xcexcm is formed, a problem arises that when the foregoing heat treatment is effected-to obtain the foregoing high piezoelectric strain constant, cracking can occur in the film. As described in JP-A-3-232755, it is considered that a bulk ceramic having a higher density exhibits better piezoelectric characteristics. However, in order to make a good application of a film having a very high density to an actuator for ink jet recording head, etc., the optimum thickness of the piezoelectric film is from about 0.5 to 25 xcexcm. When a piezoelectric film having this thickness is produced at a single step, it is normally liable to cracking. If thin films are laminated to avoid cracking, it requires a prolonged production process which is industrially unsuitable.\nFurther, an approach for raising the thickness of the piezoelectric film by repeating a process which comprises applying a sol or gel composition to a substrate, and then calcining the material at a high temperature is disclosed in xe2x80x9cPhilips J. Res.xe2x80x9d, 47 (1993), pages 263-285. However, this approach is disadvantageous in that the resulting thin piezoelectric film not only has a laminated interface that makes it impossible to provide good piezoelectric characteristics but also exhibits a deteriorated workability.\nIn general, a thin piezoelectric film is formed on a metal film which has been formed as a lower electrode on a substrate. However, a problem arises that the heat treatment effected during the formation of this thin piezoelectric film causes the substrate to be warped or distorted. Further, it is necessary that a good adhesion be established between the lower electrode and the thin piezoelectric film.\nJP-A-50-145899 discloses a piezoelectric element comprising a bulk ceramic suitable for the generation of a high voltage. However, this differs in purpose from the present invention, which concerns a thin piezoelectric film element which can be applied to an ink jet recording head.\nU.S. Pat. No. 5,265,315 discloses an ink jet recording head similarly to the present invention. However, this patent has no reference to the pores in PZT as piezoelectric film or the density thereof. Further, the proposed process for the preparation of the piezoelectric film comprises the use of sol-gel method and thus requires the lamination of a plurality of layers and a heat treatment process. Therefore, this proposal is industrially unsuitable.\nIn the above cited JP-A-6-116095, the orientation by X-ray diffraction wide angle method, i.e., orientation in the plane along the surface of the substrate is discussed. However, X-ray diffractometry of thin film is not discussed.\nFurther, if a piezoelectric element is used as an actuator for ink jet recording apparatus, etc., good piezoelectric characteristics are required. However, the relationship between crystal orientation and piezoelectric characteristics is not disclosed in JP-A-6-116095\nIt is therefore an object of the present invention to provide a thin piezoelectric film element having improved piezoelectric characteristics and a process for the preparation thereof.\nIt is another object of the present invention to provide a thin piezoelectric film element having a high piezoelectric strain constant.\nIt is other object of the present invention to provide a thin piezoelectric film element which can be prepared to have a necessary thickness without being cracked.\nIt is still other object of the present invention to provide a process for the preparation of a thin piezoelectric film element which can provide a thin piezoelectric film element comprising a thin piezoelectric film having a necessary thickness at a single process without causing cracking.\nIt is further object of the present invention to provide a piezoelectric element comprising a thin piezoelectric film having a good adhesion with a lower electrode.\nIt is further object of the present invention to provide an ink jet recording head comprising such a thin piezoelectric film element which can provide high precision printing.\nThese and other objects of the present invention will become more apparent from the following detailed description and examples.\nThe present invention can accomplish the foregoing objects and concerns a novel improved thin piezoelectric film element. The present invention provides a thin piezoelectric film element comprising a metal film formed on a substrate and a thin PZT film comprising lead zircotitanate having a third component incorporated therein formed on said metal film, wherein said thin PZT film has a rhombohedral crystalline structure which has (100) orientation of not less than 30% as determined by X-ray diffractometry of thin film, thereby enhancing the piezoelectric characteristics of said thin piezoelectric film element.\nThe foregoing orientation can be accomplished with an arrangement such that the annealing temperature of thin PZT film is from higher than 750xc2x0 C. to lower than 1,000xc2x0 C., preferably from not lower than 800xc2x0 C. to not higher than 1,000xc2x0 C., and the molar ratio of Zr/Ti is preferably from not less than 35/45 to not more than 45/35.\nMore preferably, the crystalline structure of the thin piezoelectric film element is further improved. The present invention further provides a thin piezoelectric film element comprising a piezoelectric film made of a polycrystalline substance and an upper electrode and a lower electrode arranged with the piezoelectric film interposed therebetween, wherein the crystalline constituting the piezoelectric film, i.e., the crystal grain boundary, is formed almost perpendicular to the surface of the electrodes."}
{"text": "The present invention relates to storage apparatuses for performing asynchronous remote copying between primary and secondary disk controllers. In particular, the present invention relates to techniques effective for adaptation to the transfer of update data.\nConventionally, when performing an asynchronous remote copying between a primary disk controller of a storage apparatus and a secondary disk controller thereof, by way of update data from the primary disk controller to the secondary disk controller, all write data to the primary disk controller are transferred.\nJapanese Unexamined Patent Application Publication No. 2003-202962 discloses a storage controller proposed to achieve an object of minimizing performance deterioration even at an increased inter-controller distance. To achieve the object, after a termination notification of write data is returned from a primary disk controller, the write data is directly transmitted to a secondary disk controller, and the secondary disk controller stores the received data into a nonvolatile memory, whereby data securement is performed. In addition, predetermined reference time is set to enable all write data occurring before the reference time to be secured and to enable all data occurring after the time to be abandoned."}
{"text": "The present disclosure relates to an information processing apparatus.\nIn recent years, the importance of design is growing in information processing apparatuses such as a personal computer, a mobile terminal, and an electronic book apparatus. Further, even in one model thereof, there are many combinations of color, pattern, and the like for a casing, and the variety of designs increasingly becomes important.\nFor example, in the case where a pattern is provided on a casing, laser processing is performed to respond to a variety of designs.\nFurther, regarding a casing, a high-quality texture can be given thereto, in addition to the pattern or shape thereof, by painting an inner side of a top plate formed of an acrylic transparent resin to cause a color to stand out and providing a pattern for irregular reflection to a front surface thereof.\nIn addition, there is also known a structure in which a panel constituting a part of a casing emits light. For example, a light guide body is used as the panel. In the structure, a light source such as an LED (light emitting diode) is provided inside a device and the light from the light source is propagated to the panel of the light guide body. On a rear surface of the panel, patterns such as logos are printed in white color, and light propagated along the panel of the light guide body is partially reflected on the printed portion so as to be output to the front side of the panel. As a result, there is obtained such an illumination effect that the patterns such as logos are profiled on the front surface of the panel, and the improvement in aesthetic appearance of the device can be expected (see, for example, Japanese Patent Application Laid-open No. 2004-326901 (paragraph [0029], FIG. 3; hereinafter, referred to as Patent Document 1))."}
{"text": "1. Field of the Invention\nThe invention relates to a wall cladding panel for the outside wall of a building with a solar generator and with a frame which surrounds the panel.\n2. Description of Related Art\nThese wall cladding panels are known, for example, from U.S. Pat. No. 4,223,667. Here, the frame has a metal profile which, with one leg, overlaps a glass pane outside.\nThe disadvantage of the known wall cladding panels is that the frame is not flush with the outer surface, and moreover, installation and sealing of the wall cladding panels are complex."}
{"text": "1. Technical Field\nThe present invention relates to a printing device in which color ink nozzle arrays are arranged left-right symmetrically and a printing method.\n2. Related Art\nIn cases where a number of nozzles are formed on a print head in lines, if a print head is mounted in an inclined manner with respect to its original mounting position, the landing position will be displaced from the originally assumed position. As an example of the countermeasures, what are disclosed in Japanese Unexamined Laid-open Patent Application Publication No. 2009-000836 and Japanese Unexamined Laid-open Patent Application Publication No. 2009-149064 are known.\nWhat are disclosed in the above mentioned publications are directed to a displacement in the so-called main scanning direction as a displacement of the landing position, and utilize the data of the position displaced in the main scanning direction to land at the displaced position, assuming that the landing position will be displaced in the main scanning direction."}
{"text": "Most meat sold in supermarkets is packaged in an open container or tray, overwrapped with a transparent film, weighed, and the weight printed on a label that is attached to the overwrap, in advance. Some fruits and vegetables are similarly packaged and weighed. The package is then placed in a refrigerator, showcase or on a shelf of a supermarket, where a customer may select the merchandise which he/she wishes to purchase.\nSingle roll wrapping machines for overwrapping food packages with plastic film are well known and are used in supermarket chains. A disadvantage of the present single roll wrapping machine is that they do not provide any space for a weighing device (e.g. a load cell or scale), which is a necessary adjunct of every food overwrapping operation. Another disadvantage of such machines is that it is difficult to transport the film from the roll to the wrapping location."}
{"text": "1. Field of the Invention\nThe present invention relates to a controlling method for use in a data processing device, and more particularly, to a storage controlling method for use in a processor comprising a store port for holding store data, which is transmitted from an arithmetic unit and is to be stored in a cache memory, etc.\n2. Description of the Related Art\nMainly in a super scalar processor, etc. adopting an out-of-order method, a process of a store request is performed by assigning the store request, for example, to a store port or a write buffer, which is managed by an instruction processing device and intended to temporarily hold data to be stored in a cache memory or a memory such as a main storage.\nAs conventional techniques using such a store buffer, the following documents exist.    Document 1) Japanese patent Publication No. H6(1994)-4402 “Data Processing Device”    Document 2) Japanese Patent Publication No. H10(1998)-55303 “Memory System”\nDocument 1 discloses a technique with which a write buffer for holding at least one write address and data is comprised between a central processing unit and a cache, and write data is first written to the write buffer when a store instruction is executed, and then written to the cache storage device, in a data processing device comprising the cache storage device between the central processing unit and a main storage device.\nDocument 2 discloses a memory system in which an instruction bus and a data bus are separately arranged, 4 write buffers that are interposed in parallel between a CPU and a main storage device, and do not have an address comparator are comprised, and a data write is made to a memory via a write buffer, so that the speed of the entire system is improved.\nConventionally, a data write was directly made from a write buffer or a store port to a primary cache as described above. Additionally, a dedicated write buffer was sometimes arranged for a secondary cache memory. However, a write to a primary cache was directly made from a write buffer or a store port also in this case.\nIn recent years, however, the demand for enabling out-of-order execution with much more inflight request has been rising to improve a throughput. For example, the need for increasing the number of store ports (or write buffers) has been arising. If the number of store ports is increased to improve a throughput in correspondence with such a demand, the number of store ports (or write buffers) to be processed increases, which requires a time, for example, to select from which store port data is to be stored in a cache memory. To perform such an operation on one cycle, one cycle time must be made longer, and an improvement in the throughput cannot be expected due to an increase in the number of store ports. Accordingly, a method with which out-of-order execution with much more inflight request is enabled without degrading a throughput is demanded."}
{"text": "1. The Field of the Invention\nThe invention disclosed herein is a device in which a compressive force results from the action of a straight line motion on a body having a tapered surface.\n2. The Prior Art\nThe prior art includes devices such as disclosed in U.S. Pat. Nos. 3,737,840 and 4,252,992. In '840 a deformable, star-shaped connector with a wire receiving passage therethrough is compressed inwardly around the wire by being driven into a body having a tapered opening therethrough. The device disclosed in '992 includes a plurality of jaws which are driven down a tapered passageway and thereby gripping cable ends which were inserted into the passageway."}
{"text": "The invention relates to a valve, particularly for a structural element of microfluid technology, with valve bodies having surfaces for mutual contact which are moveable relative to each other, so as to displace the contact surfaces, wherein ducts end in one of the contact surfaces for the inflow and outflow of a fluid and an indentation for the formation of a connection between the ducts is provided in the other contact surface.\nMiniaturized valves of this type which can be used for flow control in microfluid technology are known from, for example, DE 102 27 593 B4 and U.S. Pat. No. 6,748,975 B2. The contact surfaces of the valve bodies of these known valves carry out a double function. On the one hand, they seal the inlet duct relative to the outlet duct when the valve is closed. On the other hand, they ensure that the valve is sealed relative to the surroundings in any position of the valve. The latter function is of particular importance, because fluid-technological microelements, in particular flow cells, are primarily used for examining substances which burden the environment, for example, substances containing pathogenic germs.\nThe manufacture of the contact surfaces in the quality necessary for the sealing function is very difficult. In addition, numerous requirements, among them the suitability for friction pairing, limit the possibilities of material selection. Elastic synthetic materials, which are usually suitable for sealing, which however contain softeners, usually do not qualify. In the same manner, in most cases the use of lubricants for supporting the sealing function is not permissible. For example, in the case of long-term use of the valves which impairs the sealing function of the contact surfaces there is the danger of environmental contamination.\nIn addition, in accordance with the prior art, a moveable valve body for securing the sealing function relative to the surroundings/environment must always rest against the wetted sealing surface and, thus, connected tightly with a structural component which processes the fluid. This is particularly disadvantageous if the structural component processing the fluid is a disposable article. Particularly in the case of several valves per disposable, the manufacture of the moveable valve body itself and for its fluid-tight mounting, for example, at a flow cell, becomes significantly more difficult."}
{"text": "As known, a string musical instrument, such as an acoustic guitar or a violin, has a remarkable sound richness, depending on the construction features of the sound board which derive also from the fine skill of the lute maker which produces them. The acoustic power supplied directly by an instrument can be modulated by the music maker, but it cannot go beyond a certain level and in the cases of sonorisation of large indoor or outdoor spaces, it is certainly insufficient. From here the need to amplify the acoustic signal of the instrument.\nA traditional system for amplifying volume is that of placing a microphone in front of the instrument and hence amplifying the sound with conventional amplifying equipment. This solution is not always satisfactory, because it is affected by microphone and environment quality, in addition to being prone to the problem of the Larsen effect, due to the return of the acoustic signal into the microphone, precisely in situations of great amplification; moreover, the microphone is able to pick a limited dynamic of the sound which comes out of the instrument, in addition to being strongly affected by environmental acoustics and by the microphone position with respect to the instrument.\nSomewhat better results are obtained by positioning special microphones on the mouth of the sound board, below the strings. However, this solution is not fully satisfactory either.\nIn order to improve the quality of the acquired sound, freeing oneselves from the use of acoustic microphones sensitive to changes of acoustic pressure, it has also already been suggested to use special pickups (acquisition sensors), of the type found on electric guitars.\nThe pickups used in string instruments (guitars and double basses) are substantially of two types.\nA first one is of a magnetic type and comprises a series of single or double coils of metal wire, which generate a magnetic field due to the presence of permanent magnets: this type of pickup is arranged in the proximity of the metal strings of ferromagnetic material, so that the magnetic field is “disturbed” by the vibration of the metal strings and the extent of this change is acquired as electric signal (the technical term is “variable-reluctance”). This signal detects well the vibrations of the strings which move in front of the pickup, but is unable to significantly appreciate any resonance of an acoustic board, which provides the timbre of a certain instrument.\nA second one is of the piezoelectric type and comprises a piezoelectric pressure sensor: it is positioned below a bridge, for detecting the direct or resonance vibration in terms of variable pressure, consequent to the mechanical action of the strings vibrating on the support bridge. The nature of the signal acquired with this sensor, as can be guessed, is profoundly different from the one obtained with a magnetic pickup.\nIn order to acquire a signal from a piezoelectric transducer, it is normally resorted to preamplification circuits with a high input impedance, such as the schematised ones in the attached FIGS. 1A and 1B.\nA first exemplifying circuit structure (FIG. 1A) is obtained through Q1 and Q2 and defines a discrete-component amplifier stage with a high input impedance, obtained with the use of a FET (Q1) in a bootstrap configuration. The second circuit structure (FIG. 1B) is obtained by means two operational amplifiers and defines an amplifier with a high input impedance. In both cases the circuits are adequate as voltage amplifiers.\nWith known electronic configurations, no perfect adaptation of the input stages towards the piezoelectric transducers is obtained and a significant reduction of signal quality is determined in case of parallel connection of multiple piezoelectric transducers. Moreover, since the mixing between signals coming from different, jointly-used transducers is critical, in the prior art no optimisation of the logarithmic response curve of the individual level controls for each signal is provided.\nIn order to supply an amplified sound as faithful as possible to the natural one, it has also been suggested to simultaneously acquire the signal deriving from multiple different sensors, arranged on a same acoustic instrument. WO2004/023454, for example, discloses a preamplification system for a pair of sensors arranged, on an acoustic instrument, one below a bridge for detecting the direct vibrations of the strings, and one on the instrument sound box, for detecting resonance vibrations. WO2011/003148 discloses a similar system, wherein a third sensor in the shape of a microphone is also provided.\nBoth these prior-art systems, however, still have significant management problems of the different signals. As a matter of fact, in the light of the different nature of the signals produced by the various sensors, as well as the different nature of the acquired sound (direct vibration, resonance vibration, acoustic wave in the air, . . . ), problems of electronic processing of the signals in the preamplification stage exist, to then obtain the correct output signal to be amplified in the final stage, without distorsions, without spectrum limitation of the signal, and having a final sound as natural as possible. It must be observed, moreover, that the magneto-dynamic inertia sensors have limited responses in frequency, remarkable mass and high sensitivity towards external electromagnetic fields, especially those at network frequency and harmonic frequencies thereof.\nFinally, it must be considered that the reduced space available on the musical instrument also makes critical the physical configuration of the preamplifier which must be of excellent quality as well as having high immunity to electromagnetic fields of external sources.\nU.S. Pat. No. 7,304,232 discloses a joystick gain control system for the amplification of a string instrument. Two equal magnetic-type or piezoelectric-type pickups are provided, with the volume potentiometers directly connected to the pickups. Such configuration in actual fact does not allow to obtain a good response from piezoelectric-type sensors, because it would require a load resistance of a few MΩ, necessary in order not to load the sensor, which would imply an alteration of the frequency response upon the varying of the potentiometer cursor position. Moreover, a change of the equivalent resistance seen from the piezoelectric sensor would imply a change of the frequency response on the low part of the spectrum (equivalent to a sort of response of a high-pass filter at variable frequency). The volume control is arranged between the pickups and the amplification stages. The structure is that of a generic amplification system, volume adjustment of a passive type and tone control, but it cannot be understood what the inner structure and the circuitry in the active version are. The volume controls are carried out directly on the sensors, changing the response thereof, increasing the equivalent electric noise thereof due to the additional resistance of the potentiometer circuit, virtually allowing the use of sole low-impedance sensors; an analogic→digital and digital→analogic chain is furthermore accomplished, obtaining the volume and tone adjustment functions numerically. The phase inversion of either one of the two sensors is obtained through an electromechanical-type switch (a double deviator with crossed connections).\nIn the document Jarmo Landevaara: “The Science of Electric Guitars and Guitar Electronics” some amplification circuits for a piezoelectric sensor are disclosed. This prior-art solution does not offer an ideal amplifier in the application considered by the present invention yet, because it is a FET impedance-adapter, follower circuit."}
{"text": "This invention relates generally to electrical appliances used in food preparation, and more particularly, to automatic bread making appliances."}
{"text": "The present invention relates to variable picture rate coding/decoding and conversion of scanning type of a video signal from interlaced scanning to progressing scanning, for example, in highly efficient coding of moving pictures into a bitstream at a small code amount for efficient video data transfer, storage and displaying, such as, MPEG-coding with inter-picture predictive coding.\nMoving-picture coding at a low transfer bit rate, for example, 60 kbps, decimates several pictures from an incoming moving picture and encodes the remaining pictures according to need. An incoming moving picture carrying 30 frames per second (30 fps) is decimated, for instance, to 15, 10 or 5 fps. Decimation of pictures decreases the number of pictures, or frames, thus decreasing the amount of generated codes, although, motion smoothness on screen will be degraded a little bit.\nPictures under MPEG-coding are divided into three different types I-, P- and B-pictures. I-pictures (intra-coded pictures) are coded independently, entirely without reference to other pictures. P-pictures (unidirectionally predictive-coded pictures) are compressed by coding the differences the pictures and reference preceding I- or P-pictures. B-pictures (bidirectionally predictive-coded pictures) are also compressed by coding the differences the pictures and reference preceding or upcoming I- or P-pictures.\nB-pictures can be removed from a coded bitstream for changing a picture rate because B-pictures are not used as reference pictures. A bitstream of 30 fps with P-pictures for every 3 frames, for example, can be converted into a bitstream of 10 fps by removing B-picture streams only.\nA well-known variable picture rate coding performs predictive-coding for P- and B-pictures after decimation. A changed picture rate due to decimation varies a distance for prediction between a picture to be coded and a reference picture. The lower the picture rate, the longer the prediction distance, thus the amount generated codes being not so decreased.\nParticularly, in MPEG, B-picture decimation causes lowering of picture rate too much, and can not change picture rate at several stages.\nMoreover, for interlaced moving picture, decimation in unit of field lowers vertical resolution and decimation in unit of frame causes time reversal due to frame interpolation.\nDecimation in unit of field or frame for varying picture rate is thus not applicable to an interlaced moving picture.\nAlthough it is applicable to progressive moving picture, several scanning lines are decimated in reproduction by interlaced scanning.\nThis decimation processing causes redundancy in scanning lines and decoding processing. It further causes change in picture rate that depends on the amount of generated codes. Particularly, a picture rate tends to be lowered for moving picture of a big movement on screen by decimation in progress scanning, thus un-smoothness in reproduced picture being noticeable.\nMoving picture is composed of interlaced or progressive moving pictures, as discussed above. An interlaced moving picture has been decimated half the scanning lines, however, has resolution about 70% of that in progressive scanning on stationary scenes. The progressive scanning produces almost no line flicker or crawling.\nTV broadcast usually employs interlaced scanning, however, digital TV broadcast employs both interlaced and progressive scanning. Cinema films and animation carry progress moving pictures of about 24 frames per second as an interlaced scanning signal. Such an interlaced scanning signal carries 60 fields per second with the same picture for two or three successive fields.\nEncoding of moving pictures by MPEG inter-picture predictive coding requires the same bit rate for both interlaced and progressive scanning when the picture rate and the number of scanning lines is the same. Progressive scanning is, however, superior to interlaced scanning on coding efficiency because the former carries scanning lines twice the latter.\nA picture rate converted to half in progressive scanning lowers frame correlation, thus a required bit rate being lowered to 60 to 80% of the original rate. The bit rate in progressive scanning is lowered drastically compared to that in encoding in interlaced scanning at the same number of scanning lines. A picture rate is, however, half in progressive scanning, thus loosing smoothness in motion on screen a little bit when reproduced.\nA well-known scanning line conversion performs conversion of interlaced pictures into progressive pictures at the same picture rate. Progressive pictures converted to half the picture rate lowers a coding bit rate compared to interlaced pictures having the same number of scanning lines, however, being reproduced rough on screen for rapidly moving pictures."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a server and the method for recovery from a failure in one of its links (hereinafter referred to “link recovery”), and in particular, to link recovery in a server equipped with a PCI Express interface.\n2. Description of the Related Art\nAs the computer system processes more data and the processor becomes faster, the interface for interconnecting various components in the computer system is required to transfer larger amounts of data at higher speeds.\nFor some time now, the Peripheral Component Interconnect (PCI) has been used widely as an interface for interconnecting various components in the computer system. In more recent years, the serial PCI Express has been catching broader market attention because it realizes high-speed, large-capacity data transfer at low implementation costs. It is now expected to be used widely in a variety of computer systems, from personal computers (PCs) and small-scale servers for front-end use to mission critical servers for back-end use.\nMission critical servers are required to have high system availability. It is important to minimize the possibility of a system down as well as the system's down time. One of the known methods of achieving high availability in a system equipped with a PCI Express interface has been the reduced lane mode of operation, whereby, in the event of a failure in one of the links in the system, an alternative link is configured using those lanes constituting the failing link which are usable, so that the system is kept operating with reduced lanes.\nIn such an arrangement, link recovery calls for the recognition of whether any receiver is present on the PCI Express interface. Without this information, it is impossible to tell whether the link is failing even though a receiver is present on the interface or it is failing because no receiver is present on the interface; as a result, it is impossible to properly process link recovery. In addressing this problem, PCI Express employs a mechanism called receiver detection during the link training sequence, as defined in the PCI Express Specifications (refer to the Web link below), to determine the presence or absence of a receiver.\n“PCI Express Base Specification 1.1a,” PCI-SIG <URL: http://www.pcisig.com/specifications/pciexpress/base/>\nReceiver detection is a mechanism for detecting the presence or absence of a receiver on a link, whereby, after power-on of the system, the transmitter on the PCI Express interface applies a certain level of voltage to the lanes that make up the link and measures the difference in transition time to determine whether a receiver is present on the link. It takes advantage of the fact that the presence of a receiver pulls down the receiver-side of the AC coupling capacitor inserted between the transmitter and the receiver on the physical signal lines of PCI Express, thereby increasing the load capacitance as seen by the transmitter and lengthening the signal transition by as much time as needed to charge the excess capacitance. This means that the signal transition time is long if a receiver is present on the PCI Express interface, and is short otherwise. This difference is used to determine the presence or absence of a receiver on the interface.\nAs stated above, the PCI Express interface is also expected to apply to large-scale servers that are used in mission critical systems. With servers of a large form factor such as those, implementing a high-speed interface as fast as, for example, 2.5 GHz results in extended wiring lengths, which makes it difficult to keep the signal quality because of the resulting transmission losses. Therefore it becomes necessary to insert, along the PCI Express interface, a redriver such as an equalizer for compensating transmission losses. The problem, however, is that with the receiver detection mechanism, the insertion of the redriver makes it appear as if there were always a receiver present on the interface, rendering it impossible to determine whether or not an I/O extension adapter is mounted.\nAs a result, when a link failure occurs on a PCI Express interface in a server, one cannot determine whether one should carry out link recovery, and thus cannot realize high system availability."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device, and more particularly, to a method and apparatus for correcting a preferred color, which is capable of correcting a color of an input image to a color preferred by a person so as to improve image quality, and a liquid crystal display device using the same.\n2. Discussion of the Related Art\nHigh resolution and high definition of an image display device has been realized according to user's requirements. Most users determine the definition of an image on the basis of a preferred color displayed on a display device, such as a skin color, a green color or a blue color. This is because the preferred color is stored in a color storage space of a person so as to have a significant influence on color perception. Accordingly, the image display device uses a preferred color correcting method for detecting a preferred-color area from an input image and converting the detected preferred-color area into a color preferred by the user, in order to display a high-definition image preferred by a user. In the method for correcting the preferred color, the preferred-color area should be accurately detected such that other color areas are not included, and should be corrected to the color preferred by the person.\nAs a conventional preferred color correcting method, an area correcting method for deciding an input color area and a preferred-color area in an elliptical shape in a u′v′ chromaticity coordinate and mapping the input color area to the preferred-color area (“Preferred Skin Color Reproduction Based on Adaptive Affine Transform”, IEEE Transactions on Consumer Electronics, Vol. 51, No. 1, pp 191-197, 2005) and a point correcting method for setting one point of a color space as a target and positioning an input color to be close to the target (“Skin color reproduction algorithm for portrait images shown on the mobile display”, SPIE vol. 6058, pp 1-8) were reported.\nHowever, the area correcting method is disadvantageous in that contour noise occurs and luminance deteriorates because brightness is not corrected. In addition, the point correcting method is disadvantageous in that preferred color correction capability deteriorates because the contents of the input image are not considered."}
{"text": "The invention relates generally to phosphors, specifically phosphors for use in fluorescent lamps. More particularly, the invention relates to europium and manganese activated aluminate phosphors doped with rare earth elements.\nA phosphor is a luminescent material that absorbs radiation energy in a portion of the electromagnetic spectrum and emits energy in another portion of the electromagnetic spectrum. Phosphors of one important class are crystalline inorganic compounds of very high chemical purity and of controlled composition to which small quantities of other elements (called “activators”) have been added to convert them into efficient fluorescent materials. With the right combination of activators and inorganic compounds, the color of the emission can be controlled. Most useful and well-known phosphors emit radiation in the visible portion of the electromagnetic spectrum in response to excitation by electromagnetic radiation outside the visible range.\nAluminate phosphors such as barium-magnesium-aluminate (BAM) are widely used as the blue-emitting component of the phosphor blends in most fluorescent lamps intended for white light generation. Such phosphors generally have the formula (Ba, Ca, Sr)MgAl10O17. These phosphors may contain various activator ions, which impart the phosphor property. For example, a divalent europium (Eu2+) activated phosphor absorbs ultraviolet (UV) emission (i.e., exciting radiation) from the mercury plasma in a fluorescent lamp and emits blue visible light. Furthermore, a divalent manganese (Mn2+) activated BAM phosphor produces blue-green emission in fluorescent lamps.\nDespite its wide use, BAM is notorious for its shortcomings in brightness and maintenance, particularly in those applications involving exposure to high ultraviolet (UV) and vacuum ultraviolet (VUV) fluxes. These phosphors suffer from poor efficacy and lumen maintenance, specifically under high wall load conditions, which is usually found in compact fluorescent lamps (CFLs), and linear fluorescent lamps. Efficacy is the luminosity per unit of input electric power (measured in units of lumens/watt). Lumen maintenance is the ability of the phosphor to resist radiation damage over time. Because of these shortcomings, the blue BAM emission is reduced at a significantly faster rate over time than the emissions of the other color components in the blends or pixels. This results in a loss of lumens and a color shift in the overall light output.\nIt is believed that the poor efficacy and lumen maintenance are caused by UV-induced visible absorption centers, such as “color centers” and other lattice defects. Color centers are believed to be caused by lattice defects in the lattice that trap an electron or a hole, as described on pages 79-80 of K. H. Butler, Fluorescent Lamp Phosphors, Penn State University Press, 1980, incorporated herein by reference. It has been established that in many fluorescent lamp phosphors, the color centers are created by the 185 nm radiation emitted by the mercury plasma. The color centers induce absorption of the exciting radiation anywhere from the deep UV to the infrared region of the spectrum. Thus, these centers can degrade phosphor brightness by either absorbing the visible phosphor emission or by absorbing a part of the 254 nm mercury exciting radiation.\nTherefore, it would be desirable to obtain a BAM phosphor with an improved efficacy and lumen maintenance."}
{"text": "As computer-based systems become more prevalent, the quality of the interfaces through which humans interact with these systems is becoming increasingly important. One interface that is of growing popularity due to its intuitive and interactive nature is the touchscreen display. Through a touchscreen display, a user can perform a variety of tasks by contacting a region of the touchscreen with the user's finger. In order to create a more intuitive and enhanced user experience, designers often leverage user experience with physical interactions. This is generally done by reproducing some aspects of interactions with the physical world through visual, audio, and/or haptic feedback. Haptic feedback often takes the form of a mechanical vibration. There is a need for additional systems and methods to generate haptic feedback."}
{"text": "Vector network analyzers (VNA) are instruments that measure the magnitude and phase of signals as they pass through and/or are reflected from devices under test (DUTs). Typically, a DUT is connected to the VNA at connectors with short lengths of cable; however, there are some applications where the connectors of the DUT are very far away, and it is not feasible to connect the DUT connectors to the VNA with test port cables. For instance, if the DUT is coaxial cable installed in a building, one end of the cable may be hundreds of meters away from the other. The DUT could also be a radio link with transmitting and receiving antennas positioned very far apart.\nThere is a need for network analyzers capable of measuring transmission magnitude and phase (s21) through DUTs that have input and output ports very far apart."}
{"text": "The present invention generally relates to a transmission control device for vehicles and to a steering assembly for vehicles. More specifically, the present invention relates to a transmission control device for coordinating steering inputs with speed inputs and to a steering assembly providing increased reliability and space efficiency.\nWhen the driver of a typical vehicle makes a turn, the vehicle responds by changing its direction. In certain vehicles, such as the conventional tractor, the steering wheel is coupled to a steering linkage which, in turn, is coupled to the front wheels. When the driver turns the steering wheel, the front wheels pivot clockwise or counterclockwise. In one type of tractor, commonly known as a zero turn radius tractor, the rear drive wheels rotate independent of one another. The driver controls both the speed and direction of this tractor by controlling the motion of the rear drive wheels relative to one another. In this case, the steering linkage is coupled to the rear drive wheels. When the driver turns the steering wheel, this causes the rear drive wheels to rotate at different rates, which causes the tractor to turn.\nIt has been found that, in this type of tractor, when the driver is in the process of making a turn, the radius of the turn can change even though the driver holds the steering wheel in one position. This can occur when the driver increases or decreases the ground speed while making the turn. For example, if the driver turns the steering wheel to follow along a curved driveway, and at the same time, the driver presses the foot pedal, increasing the ground speed, the tractor can slightly wander away from the curved driveway. This is because the transmission of this tractor, which controls the ratio of right rear wheel speed to the left rear wheel speed, produces a change in this ratio when the driver changes the ground speed in the midst of a turn.\nIn addition to this wandering disadvantage, this type of tractor also has the disadvantage of castor front wheels which do not function like the conventional automobile-type front wheels. The castor wheels, which are not linked to the steering wheel, are free to swivel in any direction. This freedom is necessary to prevent front wheel slippage when the tractor wanders as described above.\nNot only are castor front wheels less familiar to automobile drivers, they tend to cause the front end of the tractor to wander when the tractor is traveling laterally along a slope or hillside. For these reasons, drivers must spend time to acquire the skill necessary for steering and operating this tractor in various driving conditions.\nHowever, if the transmission of a zero turn radius tractor could be adapted to prevent the wandering problem described above, the tractor could use conventional automobile-type front wheels. If a conventional automobile-type steering linkage, commonly known as an Ackerman-type steering linkage, were to be used in a zero turn radius tractor, such a steering linkage would present several disadvantages. This steering linkage includes a rack and pinion gear assembly and steering arms, both of which are positioned behind the front axle of the vehicle. This type of linkage consumes valuable space which could be occupied by other parts of the vehicle. In addition, the rack and pinion gear assembly is relatively complex which leads to several disadvantages. The gear assembly can malfunction relatively frequently, require a relatively high amount of maintenance service and is relatively expensive to manufacture.\nThe zero turn radius tractor described above also requires a reverse travel mechanism in order for the tractor to properly respond to the driver's steering inputs while traveling in reverse. This mechanism is necessary, in part, because the steering wheel and the foot pedals are separately and independently coupled to the rear drive wheels. The reverse travel mechanism is relatively complex, including a relatively high number of mechanical and electrical parts. Therefore, the reverse travel mechanism results in a relatively significant manufacturing expense and can require maintenance, service and replacement from time to time.\nTherefore, there is a need to overcome each of the disadvantages described above."}
{"text": "Most central nervous system (CNS) pathologies share a common neuroinflammatory component, which is part of disease progression, and contributes to disease escalation. Among these pathologies is Alzheimer's disease (AD), an age-related neurodegenerative disease characterized by progressive loss of memory and cognitive functions, in which accumulation of amyloid-beta (Aβ) peptide aggregates was suggested to play a key role in the inflammatory cascade within the CNS, eventually leading to neuronal damage and tissue destruction (Akiyama et al, 2000; Hardy & Selkoe, 2002; Vom Berg et al, 2012). Despite the chronic neuroinflammatory response in neurodegenerative diseases, clinical and pre-clinical studies over the past decade, investigating immunosuppression-based therapies in neurodegenerative diseases, have raised the question as to why anti-inflammatory drugs fall short (Breitner et al, 2009; Group et al, 2007; Wyss-Coray & Rogers, 2012). We provide a novel answer that overcomes the drawbacks of existing therapies of AD and similar diseases and injuries of the CNS; this method is based on our unique understanding of the role of the different components of systemic and central immune system in CNS maintenance and repair."}
{"text": "The invention generally relates to the field of money deposit and dispensing devices, particularly devices of this kind having a plurality of compartments in a housing and a closure member associated with the opening and placed under the control of a control unit enabling the closure member to be brought into an opening position only after predetermined conditions are complied with.\nIn a known device of this kind, the individual compartments for receiving bank notes are arranged one above the other in adjacent rows. Each compartment row may be closed by a closure member of sliding door type. Initially, all of the doors are completely closed. In each row associated with determined bank notes, the respectively next compartment may be opened by opening the corresponding sliding door. On the one hand, this has the drawback that any third person will immediately determine from outside how many compartments are already empty and how many are still filled with bank notes. On the other hand, those compartments of a row associated with a first kind of bank notes which have been emptied in the meantime cannot subsequently be simply used for storing bank notes of another value, as these would then be released upon the next request of bank notes of the previous value. Furthermore, the known device is very bulky and needs much space."}
{"text": "1. Technical Field\nThe present invention relates to integrated systems on silicon chip, or SoC (“System On Chip”), comprising at least one central processing unit, or CPU, on which programs can be run, a direct memory access, or DMA, controller, and a local memory.\nThe present invention relates more particularly to such systems in which successive processes, for example using algorithms, are applied to input data, typically digital audio and/or video data. Such SoCs are, for example, included in electronic appliances such as set-top-boxes, personal digital assistants (PDA), mobile phones, and so on.\n2. Description of the Related Art\nWith the dimensions of the processing cores reducing as their processing capabilities increase, one trend is to carry out a maximum of processes via software applications, the hardware components, for example the logic gates, being used only when particularly high processing performance levels, in terms of bit rate in particular, are required.\nIn practice, the use of software applications allows the use of algorithm languages with a high level of abstraction, for example “C” language, which facilitates the design step. Furthermore, errors are corrected in software applications simply by loading a new code.\nOne trade-off between the advantages of software implementations and those of hardware implementations is to combine both aspects within one and the same system, then called “firmware”, in which the system comprises on the one hand modules comprising hardware components and on the other hand software applications which are run on the CPU, the tasks to be performed by the system being divided between the hardware modules and the software applications. When executing instructions that are part of these software applications, the CPU interacts with the hardware modules, for example by sending commands to these modules. These commands use, for example, successive processes carried out by the hardware modules, for example of digital filtering, image processing, speech recognition, MPEG encoding/decoding, and other such types, on the digital data received as input to the system.\nThe hardware modules of the “firmware” type systems therefore receive data as input, carry out processes on this received data, and deliver the processed data as output.\nIt is also necessary to optimize the processing times of the software applications that the CPU applies to the data which is stored in the local memory of the system, which the CPU accesses quickly. However, in most systems on chip, it is not possible to store in local memory all the data that needs to be processed by the system on chip, for example, all the image data in the case of a system working on images. One, or even several, mass storage memories (for example, disks), hereinafter designated external memories, thus store at respective addresses in the external memories, data intended for processing by the CPU of the system on chip.\nThe function of the DMA controller of the system on chip is to transfer data from the external memory to the local memory, when this data is needed for the processing currently being carried out by the CPU. Similarly, the function of the DMA controller is to release storage resources in the local memory by transferring data from the local memory to the external memory. The data is thus transferred from memory to memory, from a source address in one memory to a destination address in another memory.\nMoreover, in some cases, the data to be processed by the hardware modules is interchanged with the hardware modules via dedicated interfaces of the processor, normally called “streaming” interfaces, each comprising an input port named SDI (Streaming Data In) and an output port named SDO (Streaming Data Out), and a register associated with each of these ports. Each of these ports normally consists of a data bus and a few synchronization signal channels.\nSuch a configuration is represented in FIG. 1. The “firmware” type system on chip SP1 comprises a CPU 1, hardware modules 2 and 3. The system SP1 has access under control of the CPU 1 to an external or local memory 6. The CPU 1 comprises a register 4 (respectively 5) linked to the hardware module 3 via the output port SDO1 (respectively the input port SDI1). The CPU 1 also comprises a register 4′ (respectively 5′) linked to the hardware module 2 via the output port SDO2 (respectively the input port SDI2).\nThe writing, controlled by the CPU 1, in the register 4 (respectively 4′), of data from the memory 6, has the effect of applying this written data to the port SDO1 (respectively SDO2). The CPU 1 can then control the writing in the register 4 (respectively 4′), of new data from the memory 6 which will, also, be applied to the port SDO1 (respectively SDO2). The reading, controlled by the CPU 1, of the register 5 (respectively 5′) has the effect of supplying the data then applied to the port SDI1 (respectively SDI2). The writing in the memory 6 of this data is then controlled by the CPU 1.\nIn some cases, one and the same port can be used for reading and writing.\nThis solution for supplying input data and recovering output data from the hardware modules is not suitable for large volumes of data because it requires intensive involvement from the CPU.\nIn other configurations, each hardware module comprises its own DMA controller. For example, the “firmware” system SP2 diagrammatically represented in FIG. 2 comprises a CPU 7 and two hardware modules 8 and 9. The CPU 7 controls access to an external memory 6 via DMA controllers. The hardware modules 8 and 9 in practice each comprise a DMA controller 8′ and 9′. The DMA controller of each hardware module can thus be used to read the data to be processed in the mass memory 6, then, once this data has been processed by the hardware module, to write the data obtained in the mass memory 6. However, this solution increases the cost of the hardware modules. It also constitutes an obstacle to the granularity of the hardware modules. For example, a hardware module comprising a DMA controller carries out two successive processes on the input data, for example a temporal interpolation and a spatial interpolation of image data. If there is a desire to divide it into two hardware modules, one carrying out the temporal interpolation and the other the spatial interpolation, in order for certain data to be able to be subjected to just one of the two processes, it is then necessary to add a DMA controller.\nThere is therefore a need for a solution to supply input data and collect output data from the hardware modules of a “firmware” type system on chip, reducing the drawbacks of the prior art."}
{"text": "Conventionally, as a sealing device, there is one described in JP-UM-A-4-93571 Publication (Patent Literature 1).\nThis sealing device is disposed between an inner race and an outer race of a rolling bearing for a wheel. This sealing device comprises a core metal member, an elastic member fixed to the core metal member, a cross-sectionally L-shaped slinger, and a garter spring. The slinger includes an axially-extending portion, and a radially-extending portion, and the elastic member includes a radial lip always sliding on the axially-extending portion, a first axial lip sliding on the radially-extending portion, and a second axial lip disposed radially inwardly of the first axial lip and sliding on the radially-extending portion. The second axial lip has an annular groove formed in a radially-outward surface thereof.\nThe garter spring is fitted in the annular groove of the second axial lip. The garter spring presses the second axial lip radially inwardly.\nThis sealing device is formed such that in a condition in which the second axial lip is not worn, the second axial lip is not in contact with the axially-extending portion, while when the second axial lip is worn, so that a press-contacting force of the second axial lip for the radially-extending portion becomes less than a predetermined force, part of that portion of the second axial lip opposed to the axially-extending portion is brought into contact with the axially-extending portion, so that part of the above opposed portion forms a radial seal.\nThe second axial lip, when not in a worn condition, functions as an axial seal, while in a worn condition of the second axial lip, part of the above opposed portion of the second axial lip functions as the radial seal, so that this sealing device can continuously maintain a stable sealing function.\nPatent Literature 1: JP-UM-A-4-93571 Publication (FIG. 1)."}
{"text": "The present invention relates to a composite flexible frozen confection containing a distinct block(s) of a gel component, that is combined with one or more frozen dessert components, in such a manner that the composition may exhibit hand-held flexibility without significant separation of one or more of the components. The present invention further teaches a method of preparing such a food composition for frozen dessert applications.\nThe key features of the invention are the fun, or play factor and the absence of the messiness of eat that would otherwise ensue. Other aspects of multiple components are also present such as having more than one visual appearance such as color and clarity, more than one texture, mouthfeel, flavor, flavor release etc. upon consumption. In particular, the frozen dessert component provides the refreshment aspect that complements the physical strength necessary for the wobbly gel component.\nWO 99/38386 relates to a water ice containing stabilizers and having a channeled structure of air passages. It does not teach flexibility or combinations having blocks of gel components with other blocks of frozen dessert component.\nEP 0864256 teaches a way of molding an ice confection using multipart molds and liquid nitrogen as a cryogen. It does not teach a way to achieve flexibility of multicomponent products.\nJP 2000 004793 relates to an iced dessert with a jelly-like solid item coated with ice cream. Product is made by coating solid or fluid edible material with ice cream, or using concentric nozzles. It teaches viscosity control and use of sweet potato fiber as additive to ice cream and does not teach flexibility.\nJP 1999 346659 relates to a swirl design food based upon molding and nozzle devices. This does not relate to a gel, or to a method of making a gel, and does not teach flexibility.\nEP 0560052 relates to the use of a gelatin coating (not a polyanionic) upon ice cream. The teaching is for low calorie products. There is no teaching regarding wobbliness or prevention of disengagement of gel layer.\nU.S. Pat. No. 3,752,678 involves dipping an ice cream into a thixotropic batch containing alginate. This product contains the separate gel phase as a distinct component block (a coating) and emphasis is upon achieving a jelly coating, not upon wobbliness.\nIn any situation of bending a typical flexible material (like a piece of eraser), it is not just the change in the direction of the arc that must be considered. There is typically one surface of the flexible material that suffers dilation (the outside of the arc) and the opposite surface suffers compression (the inside of the arc). Both dilation and compression effects typically lead to other indirect but significant changes in surface geometry. These indirect changes have large contributions to the disengagement of the components of the composition. The present invention minimizes both the occurrence and the consequences of these side effects.\nThese side effects are explained as follows.\nA dilated surface also typically undergoes some narrowing in one direction, to compensate for the stretching in the other direction. A compressed surface also undergoes some lateral broadening and buckling, in order to tolerate the material displacement that is resultant from the squeezing of its fabric in the other direction. Although these changes are side effects of the creation of the bend, they are nonetheless substantive, geometric displacements.\nWhen such changes in surface geometry of a flexible gel occur in juxtaposition with a more rigid conjoined structure such as a frozen dessert; surface separation and delamination of the union is wholly expected. Bending movements lead to substantial breaking of at least one of the more brittle non-gel components. Such broken fragments of the non-gel component(s) then typically fall off the gel component.\nThus there is a need for a composition in which two or more components of different rigidity are combined, yet in which the application of bending motions does not cause the falling apart of the composition in any significant way. Another need is for a composition that, despite being able to be wiggled about, does not generate a very messy eating experience. The present invention now satisfies these needs.\nSurprisingly, it has been discovered that with the present composition, compositions having greater structural integrity are provided such that meesy eating experiences due to breakage are substantially minimized. In fact although the more rigid component does still crack, or craze to some degree, this is not excessive. Further, the more rigid component does not readily separate or disengage from the gel component and survives repeated wiggles without disengaging. Part of this is due to the discovery that the adhesiveness of the gel phase should not be too high.\nThe invention thus concerns a composite frozen confection containing one or more component block(s) of polyanionic gel(s) in conjunction with more rigid component such that:\na) During warming to room temperature the confection develops flexibility,\nb) When flexible, the confection may be bent or wobbled such that one end may distort to an angle of at least 5 degrees from other end and\nc) During the first 4 distortions, not more than 5% by weight of the rigid component falls off.\nIn this invention the gel phase is separate, as a distinct and elongated component block, such as a filling, a layer, a rope, a coating, or as laminations etc., provided the gel phase substantially extends longitudinally from one end of the product to the other end.\nThe gel component is not present as an emulsion phase or dispersion (in which individual particulate gel sizes and geometries are not visible to the eye of a healthy individual).\nThe adjunct frozen dessert component is typically a water ice, but may be a sorbet, sherbet, ice cream etc. and might or might not itself also be, or contain, a gel. There is a difference in the rigidities of the two components of the present composition.\nThe method for preparing such items is embodying the ability to form a viscoelastic charged gel component and a rigid component together into a composite.\nThe method comprises contacting the gel component with the rigid component by processes such as molding, enrobing, spraying or co-extrusion."}
{"text": "1. Field of the Invention\nThe present invention relates to a position detection and anti-pinch system for a power window of a vehicle and more particularly to a position detection and anti-pinch system and method for detecting the position of a window in a vehicle door assembly.\n2. Background Art\nDoor assemblies for automotive vehicles generally include a metal door frame attached to the body of a vehicle. The door frame generally includes inner and outer door frame panels joined about a common edge and a door sill running along top edge of the door to define a window frame. The window frame includes an opening or channel extending through the door frame to receive a window pane. A reinforcement panel may be attached to the inner door panel to provide additional structural support for the door frame and to receive various hardware components. Hardware components such as a handle assembly, window regulator and motor assembly, and electrical components such as interior lights and speakers, are secured to the reinforcement panel concurrently with assembly of the vehicle door.\nCertain vehicle door assemblies include automatic vehicle window adjustment having a “one-touch up” feature. The “one-touch up” feature allows a passenger to raise the window pane with a single touch of a control button on the vehicle door. When actuated, the switch activates a motor assembly operatively connected to a window regulator, causing the window regulator to move the window pane upward toward the top of the window frame to a fully closed position without additional operator intervention. Of particular concern with such “one-touch” systems is that the passenger may not be able to stop the upward travel of the window pane if an object or an appendage of a person becomes pinched or trapped between the moving window and the window frame, which may lead to significant bodily harm.\nSeveral types of obstruction detection systems have been proposed for vehicle door assemblies. One known obstruction detection system includes a control module which senses current spikes in the window drive motor or pressure on the window which would indicate an obstruction being pinched between the upward moving window and the window frame. The obstruction detection system stops and/or reverses the motor upon sensing the current spike or increased pressure while raising the motor. The problem with this design is that an object must be trapped or pinched with sufficient force for the system to sense the need to reverse the motor and lower the window, causing undesirable injury to an occupant prior to the lowering of the window. Further, the sensor is mounted adjacent the motor rather than the window pane, which does not provide an accurate measurement if an object becomes trapped between the window and window pane.\nAnother known obstruction detection or anti-pinch system includes a window seal with a gasket having varying resistivity based upon its level of compression. For example, if an obstruction is pinched between the window and the window frame, the gasket is locally compressed, which alters the resistivity of the gasket. A low current is applied through the gasket, and the variation in resistivity may be sensed to determine the existence of the obstruction, at which point upward movement of the window would be stopped. The system is not particularly desirable because the resistive gasket and associated electronics may be costly and injury may still occur prior to the lowering of the window.\nGenerally, it is known in window control designs to determine the velocity of a power window via pulses from the rotation of a motor's shaft. These pulses are generated by Hall sensors (magnetic sensing), CPD (commutator pulse detection) or switches that provide an open or closed circuit as the motor spins. To coordinate a system to detect information about a moving power windowpane, a supplier must change the system that they are designing to accept these pulses. Control algorithms must then be used to determine the velocity and position of the windowpane based on the received the pulse information. This task is especially critical when considering safety standards like FMVSS 118, where an obstacle must be detected during an auto-up function before a predetermined force is exerted on the window as a result of an obstruction in the path of the power window.\nWith such stringent guidelines, the control system in such motor pulse generating arrangements must “learn” the characteristics of not only the motor, but the regulator, grease and all of the mechanical linkages between the motor and the window. The system must also be able to accommodate temperature variation and different road conditions into the control process. All of these factors must be considered to provide flawless auto-up operations.\nAccordingly, it is desirable to provide an improved window position detection and anti-pinch system for an automatic power window of a vehicle. It is also desirable to provide a window position detection system incorporating an anti-pinch detection system which prevents an obstruction from becoming trapped or lodged between a window and a window frame of the vehicle door."}
{"text": "1. Field of Invention\nThe present invention relates to a snapshot mechanism. More particularly, the present invention relates to a snapshot mechanism adopting a rollforward and a many-to-one management method, which is adaptive to a remote snapshot mechanism.\n2. Description of Related Art\nTo backup data for security, selection of a backup technique should take possible disasters, such as natural disasters, physical damage to the storage hardware and data corruption (including human errors, software errors, viruses, hacker invasion) into consideration. Only generating copies of the original data cannot solve the problem of data corruption. Thus, to correctly recover the data, the concept of version control should be adopted, and a snapshot technique meets the requirement.\nThe first advantage of the snapshot technique is that a current service will not be interrupted, by which the service level of the company will be improved. The second advantage is that a backup window is shortened, so frequently conducting snapshoot and backup the important data can be easily achieved. As such, a data loss window can be effectively shortened to reduce the estimated amount of data loss when disasters occur. The third advantage is that snapshot can solve the problem of data corruption. The above-mentioned backup window refers to the time needed for executing backup. The backup window is usually defined by the time needed for operating a backup process. For example, if data is needed from 8 a.m. to midnight, the window can be used to fabricate backup copies from midnight to 8 a.m. of the next day. However, to obtain consistent backup, data cannot be varied when being backed up. Therefore, in some circumstances, the backup window is also a time interval in which data and application programs cannot be used.\nCopy-on-write (CoW) is the mainstream of practical snapshot manners, as it has good flexibility, scalability and can be practiced at a block level. Besides, the efficiency of snapshot at a block level is higher than at a file system level. As the CoW is independent from the file system in design, the dependency thereof to the file system does not need to be considered. In the snapshot of copy-on-write, when CoW occurs, the write action is suspended till the CoW action is finished. As for a remote snapshot, the time needed to suspend a write process will add double transmission time, the application performance is affected and weakened significantly.\nIn the development of snapshot techniques, LSI Logic Corporation Company proposed a snapshot technique in the U.S. Pat. No. 6,771,843 titled “Data Timeline Management Using Snapshot Volumes” published on Aug. 3, 2004. As shown in FIG. 1, the data loss window is shortened to reduce loss and avoid accidents caused by man-made improper operations. The snapshot technique of rollforward is achieved by reserving the snapshot data after the time point for rollback plus timeline management skills.\nAs shown in FIG. 1, according to the technique, snapshot volumes 110, 120, 130, 140 respectively reserve data at one o'clock, two o'clock, three o'clock and four o'clock on one day afternoon. When rolled back to the status reserved by the snapshot volume 120, the data is checked sequentially by an algorithm provided by the LSI Company, and the data required to be reserved is replicated into one or several appropriate volumes. Thus, to roll forward, the snapshot volume 130 or the snapshot volume 140 can be used to roll forward to the status at three o'clock or four o'clock. When system administrator makes an improper operation, the system can be recovered to the point-in-time before improper operation happened. To rollback to the snapshot volume 120, for example, the data of the snapshot volume 120 will be cleared, and Copy-on-Write (CoW) must be performed again for rolling forward, thus prolonging the recovery window in such a situation. Another disadvantage is that the block of the CoW may need to be replicated into a plurality of snapshot images, which results in the waste of resources.\nIn addition, LSI Logic Corporation Company proposed another snapshot technique in the U.S. Pat. No. 6,594,744 titled “Managing a snapshot volume or one or more checkpoint volumes with multiple point-in-time images in a single repository” published on Jul. 15, 2003, which supports a single snapshot volume to store a plurality of snapshot images. As shown in FIG. 2, all the snapshots belonging to the same target volume are stored in a big snapshot volume. For example, each snapshot of a target volume 210 is respectively stored in a snapshot volume 220 in sequence, for example, snapshots A, B, C and D are respectively stored in images A, B, C and D of the same snapshot volume 220. Similarly, each snapshot of the target volume 230 is respectively stored in a snapshot volume 240 in sequence, for example, snapshots E, F, G and H are respectively stored in images E, F, G and H in the same snapshot volume 240. That is, different target volumes respectively have the corresponding snapshot volumes."}
{"text": "In a manner which in itself is known, unlike an electrical machine with permanent magnets, a rotary electrical machine with excitation can produce engine torque, or supply electric energy, only when an excitation current passes through its inductor.\nA common type of rotary electrical machine with excitation, which is widely used in the motor vehicle industry for the alternator and starter functions, comprises a rotary inductor and a stator with a plurality of windings.\nWhen the machine is functioning as an alternator, the current which is generated in the stator windings by the rotating inductor is rectified such as to supply direct voltage to the vehicle battery.\nThis voltage depends on the speed of rotation of the inductor, the connected load and the excitation current.\nFor motor vehicle applications, the output voltage must be regulated such as to remain constant irrespective of the speed of rotation of the alternator and irrespective of the battery load.\nFor this purpose, the output voltage is measured and compared continually with a set value by a control device which controls the excitation current such as to cancel out any difference.\nThe company VALEO EQUIPEMENTS ELECTRIQUES MOTEUR has already proposed to carry out this control on the basis of measurements by sampling by means of digital techniques, which provide substantial advantages in comparison with the conventional analogue methods, in particular in its European patents EP 0 481 862 and EP 0 802.606.\nIn the design of a modern control device, the subjection of the output voltage to a set value is based on the theorisation of a proportional (P) or proportional integral (PI) control loop.\nThe creation of this loop by the corresponding algorithms makes it possible to design regulators with programmable functionalities which can adapt more easily to the specifications of the motor vehicle manufacturers, such as the one described in the article “An high-voltage CMOS voltage regulator for automotive alternators with programmable functionalities and full reverse polarity capability”, P. Chassard, L. Labiste, P. Tisserand et al, Design, Automation and Test in Europe Conference and Exhibition (DATE), 2010, EDAA.\nIn the motor vehicle industry, a plurality of characteristics of the alternator make it possible to evaluate the performance of the alternator, in particular the following characteristics:\nmaximum current supplied according to the speed of rotation for a given regulation voltage;\nvoltage regulation, i.e. the aptitude of the alternator to generate a voltage corresponding to the set value according to the loads connected;\nstability of the system, i.e. the phase margin and the gain margin of the system comprising the regulator, the alternator, a battery and loads.\nHowever the inventive body has observed a decrease in the phase and gain margin in systems comprising high-power alternators.\nThis decrease can become critical. In fact poor voltage regulation arises which leads to an “oscillating” voltage, or worse still, an alternator output voltage which can no longer be controlled by the regulator. This lack of control can lead for example to excess voltage; reference is then made to an unstable controlled system.\nAccording to arbitrary criteria, the phase margin of the transfer function in an open loop (FTBO) of a controlled system must be more than 45°, and the gain margin must be more than 13 dB in order to consider the system as stable.\nHowever, mostly, for a system which uses a very high-power alternator (which supplies a current of 300 A for example), the respective criteria of the phase margin (>45°) and gain margin (>13 dB) can no longer be complied with for a performance level of the voltage regulation which is identical to that of an alternator with lower power (which for example supplies a current of 100 A). The system is liable to be unstable.\nA solution which is well-known in order to compensate for the increase in alternator gain is to create a decrease in the regulator gain.\nUnfortunately, in the case of a control loop of the proportional type, the performance of the voltage regulation is affected by the change of the regulator gain.\nIn fact, although the criteria of stability are improved, the drop in voltage of the voltage regulation according to the current output is then increased, and can reach 600 mV, which is considered to be a deterioration of the performance of the voltage regulation in comparison with a voltage drop of approximately 200 mV for an alternator which outputs 100 A, for example.\nA known solution for limiting the voltage drop of the regulation when the regulator gain is decreased is the use of an integral part in the control loop.\nThe voltage drop can thus be brought to a value close to 0 mV; however, it is known that the phase margin and the gain margin are slightly affected.\nThere is therefore a need for a solution which would make it possible to increase the phase margin and the gain margin of the transfer function in an open loop of a voltage regulation system associated with a high-power alternator, whilst maintaining the performance of the voltage regulation."}
{"text": "Clustered Regularly Interspaced Short Palindromic Repeats (CRISPR) and CRISPR-associated (Cas) systems are prokaryotic immune system first discovered by Ishino in E. coli. Ishino et al. 1987 (Journal of Bacteriology 169 (12): 5429-5433(1987)). This immune system provides immunity against viruses and plasmids by targeting the nucleic acids of the viruses and plasmids in a sequence-specific manner.\nThere are two main stages involved in this immune system, the first is acquisition and the second is interference. The first stage involves cutting the genome of invading viruses and plasmids and integrating segments of this into the CRISPR locus of the organism. The segments that are integrated into the genome are known as protospacers and help in protecting the organism from subsequent attack by the same virus or plasmid. The second stage involves attacking an invading virus or plasmid. This stage relies upon the protospacers being transcribed to RNA, this RNA, following some processing, then hybridizing with a complementary sequence in the DNA of an invading virus or plasmid while also associating with a protein, or protein complex that effectively cleaves the DNA.\nThere are several different CRISPR/Cas systems and the nomenclature and classification of these has changed as the systems are further characterized. In Type II systems there are two strands of RNA, a CRISPR RNA (crRNA) and a transactivating CRISPR RNA (tracrRNA) that are part of the CRISPR/Cas system. The tracrRNA hybridizes to a complementary region of pre-crRNA causing maturation of the pre-crRNA to crRNA. The duplex formed by the tracrRNA and crRNA is recognized by, and associates with a protein, Cas9, which is directed to a target nucleic acid by a sequence of the crRNA that is complementary to, and hybridizes with, a sequence in the target nucleic acid. It has been demonstrated that these minimal components of the RNA-based immune system could be reprogrammed to target DNA in a site-specific manner by using a single protein and two RNA guide sequences or a single RNA molecule. The CRISPR/Cas system is superior to other methods of genome editing involving endonucleases, meganucleases, zinc finger nucleases, and transcription activator-like effector nucleases (TALENs), which may require de novo protein engineering for every new target locus.\nBeing a RNA-guided system, CRISPR/Cas systems can be prone to issues with RNA-DNA hybrid structures, such as RNase A degradation of the RNA strand and higher possibility of RNA-DNA mismatches. Furthermore, synthesis of DNA oligonucleotides is more economical and robust than synthesis of RNA oligonucleotides. DNA-guided CRISPR systems may also recruit additional machinery to a specific target, compared to naturally occurring RNA-guided CRISPR systems. A need exists for an improved system that overcomes the problems associated with RNA based CRISPR/Cas systems, provides access to the decreased cost and increased robustness of DNA synthesis, and improves the specificity of the CRISPR/Cas system."}
{"text": "1. Field\nMicro-electro-mechanical (MEMS) devices and systems and in particular, but not exclusively, MEMS probe based memory storage.\n2. Background\nSolid state memories typically employ micro-electronic circuit elements for each memory bit. Since one or more electronic circuit elements are required for each memory bit (e.g., one to four transistors per bit), these devices can consume considerable chip “real estate” to store a bit of information, which limits the density of a memory chip. The primary memory element in these devices is typically a floating gate field effect transistor device that holds a charge on the gate of field effect transistor to store each memory bit. Typical memory applications include dynamic random access memory (DRAM), static random access memory (SRAM), erasable programmable read only memory (EPROM), and electrically erasable programmable read only memory (EEPROM).\nPhase-change memory uses a phase-change material as the data storage mechanism and offers significant advantages in both cost and performance over conventional memories based on charge storage. Phase-change memories use phase-change materials that can be electrically switched between two or more phases having different electrical characteristics such as resistance. One type of memory element, for example, uses a phase-change material that can be electrically switched between a generally amorphous phase and a generally crystalline local order, or between different detectable phases of local order across the entire spectrum between completely amorphous and completely crystalline phases.\nThe phase-change memory can be written to, and read from, by applying current pulses that have the appropriate magnitude and duration and that cause the needed voltages across and current through the volume of phase-change material. A selected cell in a phase-change memory can be programmed into a selected state by raising a cell voltage and a cell current for the selected cell to programming threshold levels that are characteristic of the phase-change material. The voltage and current are then typically lowered to quiescent levels (e.g. essentially zero voltage and current) that are below the programming threshold levels of the phase-change material. This process can be performed by the application of, for example, a reset pulse and a set pulse which can program the cell into two different logic states. In both of these pulses, the cell voltage and cell current are caused to rise at least as high as certain threshold voltage and current levels needed to program the cell. Next, to read the programmed cell, a read pulse can be applied to measure the relative resistance of the cell material, without changing its phase. Thus, the read pulse typically provides a much smaller magnitude of cell current and cell voltage than either the reset pulse or the set pulse.\nElectrical memory devices employing phase-change material typically do not use field effect transistor devices, but comprise, in the electrical context, a monolithic body of thin film chalcogenide material. As a result, very little chip real estate is required to store a bit of information, thereby providing for inherently high density memory chips. The phase-change materials are also truly non-volatile in that, when set in either a crystalline, semi-crystalline, amorphous, or semi-amorphous phase representing a resistance value, that value is retained until reset as that value represents a physical phase of the material (e.g., crystalline or amorphous)."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nDuring various procedures, two or more components can be combined to form a mixture of the materials for various purposes. For example, a first and second material may be mixed to polymerize, seal, or the like. For example, thrombin and fibrinogen can be mixed to form a tissue sealant that is biocompatible.\nAutologous material can be used to form a tissue sealant. The tissue sealant materials, including thrombin, can be extracted or concentrated from a patient, such as from a whole blood sample. The tissue sealing materials can be provided or held separately until applied to an area to be sealed.\nA tissue sealant can be used in various procedures, such as a minimally invasive procedure to speed healing. In addition, an autologous material can provide various benefits such as reducing or eliminating cross contamination or infection. The materials that can form a tissue sealant, however, are generally held separately until applied to a selected area."}
{"text": "Electrohydraulic valves (EHVs) operated by a torque motor are widely used in industrial applications. For example, single-stage electrohydraulic valves may be used in myriad systems and environments. One typical system and environment is the fuel control system on-board a jet-powered aircraft. No matter the specific end-use system and environment, single-stage EHVs typically include at least one nozzle that is disposed between a pressurized hydraulic fluid source and a hydraulically controlled load. Pressurized hydraulic fluid flow through the nozzle, and thus to the hydraulically controlled load, may be controlled via the torque motor.\nA conventional torque motor that is used with a single-stage EHV includes a plurality of coils, an armature assembly, and a flapper. The coils are controllably energized to control the rotational position of the armature assembly. The flapper is coupled to the armature assembly and extends between the outlet of the nozzle and a flapper stop in the single-stage EHV, defining a “nozzle-flapper assembly.” The flapper in the nozzle-flapper assembly is conventionally a steel bar. By controlling the rotational position of the armature assembly, the position of the flapper relative to the nozzle outlet in the nozzle-flapper assembly is controlled and thus fluid pressure and/or flow to the hydraulically controlled device is controlled. When the EHV is used as a shut-off valve, the flapper is held against the nozzle outlet in a closed position to prevent flow of the pressurized hydraulic fluid therefrom to the hydraulically controlled device. When the flapper is off the nozzle outlet, the EHV is in an open position, permitting flow of the pressurized hydraulic flow.\nConventional nozzle-flapper assemblies exhibit limited life and excessive leakage between the flapper and the nozzle outlets even in the closed position due to wear of the flapper and/or nozzles in high vibration environments. Over time, the leakage increases as the wear increases.\nHence, there is a need for flapper assemblies for torque motors of electrohydraulic valves. In particular, there is a need for flapper assemblies that may be used in nozzle-flapper assemblies of torque motors of electrohydraulic valves, that exhibit relatively less leakage and extended operating life, and that are relatively inexpensive to manufacture. The present invention addresses at least these needs."}
{"text": "1. Field of the Invention\nThis invention relates to a toner cartridge which is detachably mounted on a toner box of a copier/duplicator or a printer to supply toner into the latter.\n2. Description of the Related Art\nThere are available a variety of toner cartridges that are detachably mounted on the toner box. Toner cartridges, that are mounted on the toner box after opening are well known in the art. However, a toner box of this type is disadvantageous in that in the process of mounting it on the toner box, the toner is often scattered, thus contaminating the surrounding area. In order to eliminate this difficulty, toner cartridges have teen proposed, for instance, by Japanese Patent Application (OPI) No. 224364/1983 (the term \"OPI\" as used herein means an unexamined published application\"), and Japanese Utility Model Application (OPI) No. 77151/1983, where the cartridges are opened after being mounted on the toner box. Another toner cartridge is proposed by this inventor in U.S. Pat. No. 4,834,246.\nOne example of a conventional toner cartridge is shown in FIGS. 8, 9 and 10. FIG. 8 is a sectional view of the cartridge, FIG. 9 is a sectional view taken along line A--A of FIG. 8, and FIG. 10 provides a bottom view for a description of the opening of the conventional toner cartridge. The conventional toner cartridge comprises a toner container body 1 that has an opening at the bottom and is substantially in the form of a rectangular box. A mouth ring member 2 has a bonding region 8 along the opening edge to which a sealing film 4 is bonded. The mouth ring member 2 also holds a lid 3 in such a manner that the lid is slidable. The sealing film seals the opening in such a manner that it can be peeled off when necessary. At one nd of the opening, a locking member 11 is present to secure one end of the sealing film and the lid for the opening. Toner 6 is contained in the toner container body, which is covered with a top plate 7. The sealing film is bonded to the bonding region of the mouth ring member, thus supporting the toner. One end portion of the sealing film, is folded over one end portion of the lid, and secured to the locking member 11. The locking member 11 is slidably provided over the rear surface of the lid and the lid is slidably held by a side wall 10, which is integral with the mouth ring member.\nThe toner is supplied by the conventional toner cartridge as follows. First, the toner cartridge is mounted or the toner inlet of the toner box. Then, the locking member is moved in the direction of the arrow of FIG. 8 with a suitable means (not shown) to peel the sealing film off the opening of the toner container body thereby opening the latter. As a result, toner is supplied to the toner box from the toner container body.\nIn the conventional toner cartridge described above, the bonding region of the mouth ring member is adjacent to the opening. If the toner cartridge is dropped during transportation or handling, it often occurs that the sealing film falls off, or is damaged or punctured by the impact of dropping, resulting in toner leaking out of the toner cartridge. This difficulty may be eliminated by strongly bonding the sealing film to the bonding region but this results in requiring greater power to peel off the sealing film to supply the toner, resulting in a toner supplying operation that cannot be easily or smoothly carried out.\nAccordingly, an object of the present invention is to provide a toner cartridge that has the seemingly contradictory characteristics of allowing the peeling of the sealing film to be easily and smoothly performed while providing a toner cartridge configuration that reduces damage to the sealing film as a result of accidental dropping of the toner cartridge during handling or transportation.\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out in the appended claims."}
{"text": "Polar organisms should overcome the problems of decreased enzyme activity, decreased membrane fluidity, inactivation and improper folding of proteins, formation of intracellular ice crystals, etc. to survive in low-temperature, polar environments. Among others, the formation of ice crystals causes physical damages and dehydration of tissues due to the growth of ice crystals, thus causing serious damage to polar organisms. Polar organisms produce various antifreeze proteins (hereinafter referred to as “AFPs”) to survive at low temperatures. AFPs inhibit the growth of ice crystals in vivo and the recrystallization of ice to protect polar organisms from sub-zero temperatures to survive (Davies, P. L. and Sykes, B. D., Curr. Opin. Struct. Biol. 7, 1997, 828-834; Davies, P. L. et al., Philos Trans R Soc Lond B Biol Sci. 357, 2002, 927-935; D'Amico, S. et al., EMBO Rep. 7, 2006, 385-389).\nAFPs are proteins that generally have a flat ice-binding surface and bind to specific surfaces of ice crystals, thus inhibiting the growth of ice crystals and the recrystallization of ice. AFPs create a difference between the melting point and freezing point. This is called thermal hysteresis (TH), which can be measured using a nanoliter osmometer and used as an indicator of AFP activity. Moreover, AFPs do not lower the freezing point in proportion to the concentration, unlike typical antifreeze used in vehicles. That is, AFPs can effectively lower the freezing point even at very low concentrations by direct interaction with ice, thus minimizing damage due to osmotic pressure generated in vivo during freezing (Jia, Z. and Davies P. L., Trends Biochem. Sci. 27, 2002, 101-106).\nThe unique features of AFPs that prevent the growth of ice crystals and inhibit the recrystallization of ice have been used in various commercial fields. For example, in the agricultural field, AFP expression in plants has been attempted for the purpose of preventing cold-weather damage to plants. Moreover, in the field of fisheries, there has been an attempt to produce a transgenic fish by expressing AFPs in commercially available fish such as Atlantic salmon (Salmo salar) or goldfish (Carassius auratus) so as to enable farming in cold areas. Furthermore, in the medical field, research on the use of AFPs in cryosurgery and as an additive in cryopreservation of blood, stem cells, umbilical cord blood, organs, and germ cells has continued to progress. In addition, in the food field, AFPs are also used in product production for frozen storage of smoother ice scream. In the field of cosmetics, functional cosmetics containing AFPs for preventing frostbite have already been sold. Although AFPs are widely used in various commercial fields as mentioned above, there are still limitations in mass production of recombinant AFPs due to low-level expression of AFPs and folding problems. This is mainly because most AFPs have disulfide bonds and are stabilized by disulfide bonds, which thus makes it difficult to express recombinant proteins and yields improper folding of expressed proteins.\nSince AFPs were first discovered in fish living in cold water, various types of new AFPs have been discovered in insects, plants, fungi, microorganisms, etc. New AY30 AFP derived from arctic yeast, Leucosporidium sp., has recently been recovered. The AY30 AFP has no cysteine amino acid residues, and thus during production of recombinant proteins, the level of protein expression is high, and the folding problem due to improperly formed disulfide bonds does not occur, As a result, the AY30 AFP is suitable for mass production of recombinant AFPs.\nTherefore, the present inventors have synthesized a recombinant polynucleotide by modifying an AFP gene to be expressed using codon optimization for a yeast expression system and inserted the recombinant polynucleotide into a yeast-derived expression vector so as to mass-produce an antifreeze protein (AFP) derived from arctic yeast by overexpressing AFP in the form of activated protein. As a result, the present inventors have obtained a large amount of AFP and found that the AFP is glycosylated, thus completing the present invention. All references cited in this specification are hereby incorporated by reference in their entirety."}
{"text": "1. Field of the Invention\nThe invention relates to a method for mobile on-line and off-line monitoring of colored and high-gloss automobile component surfaces. For evaluating the surface quality of series-produced coated and/or painted automobile body components, it is standard procedure at the present time to analyze quality parameters such as color, gloss, layer thickness and wave (waviness) to determine different characteristics, dependencies and parameters. Commercially available manual measuring devices are used for this, for example the angle spectrometer X-Rite MA 68 II and gloss-measuring devices such a Wave Scan plus by the company BYK-Gardner GmbH, 82534 Geretsried.\nFor measuring the colored and high-gloss automobile component surfaces, coating samples—so-called test panels—are generally produced which are then used to determine the optical characteristics to be examined, such as color, gloss, layer thickness and wave of the coating, as well as to analyze the mutual dependencies of the coating characteristics.\nAs a rule, this is done in the laboratory or during the production by taking random samples using the aforementioned measuring methods and devices.\n2. Related Art\nReference DE 19709406 A1 discloses a method and a device for measuring painted test panels to determine the surface quality—color, gloss, layer thickness and wave—with the aid of a laboratory robot in combination with a corresponding measuring method.\nThe disadvantage of this method and device is that the evaluation of the measuring results with the aid of test panels does not take into account the geometric form of automobile component surfaces, thus allowing only indirect conclusions to be drawn for the quality analysis.\nReference DE 19717593 A1 discloses a measuring method for evaluating the surface quality of motor-vehicle bodies by detecting in a non-contacting manner the surfaces of series-produced, automatically coated motor vehicle components in conjunction with a multi-axis robot moving along pre-programmed movement paths.\nThis invention has the disadvantage of high investment costs for the robot required for use and the associated program-technical links for controlling its measuring movement along the automobile component surface to be measured, as well as the stationary connection.\nAlso known from prior art are non-contacting measuring techniques using cameras, which are aimed at defined angles onto the automobile component surfaces to be measured and which measure these surfaces under various types of lighting and lighting angles.\nHowever, these measuring techniques do not meet the requirements of the automobile industry with respect to mobile use, degree of automation, low investment costs, flexibility of the measuring requirements, ability to simultaneously measure various measuring variables such as color, gloss, layer thickness and wave and the option of a data-technical analysis. Furthermore, a variable measuring of critical automobile component surface parameters on-line and off-line is not possible with the presently used methods, but is desired by the automobile industry.\nIt is the object of the present invention to develop a method that allows a mobile on-line and off-line monitoring of the quality of colored and high-gloss automobile component surfaces, so that parallel measurements of the parameters for color, gloss, layer thickness and wave can be realized quickly and structured according to different requirements. The core of the invention is the mobile on-line or off-line measuring through optical scanning of the colored and high-gloss automobile component surfaces with an angle-dependent spectrophotometer during the production or final control.\nThe measuring beam, formed with polarized light of different wave lengths, is thus the measuring beam for the angle-dependent spectrophotometer which is combined with a reference beam for the angle-dependent spectrophotometer and contains the reflection, interference, depolarization and phase values of the measured automobile component surface for different wavelengths as surface information, wherein these represent locally precise images of the optical surface conditions of the automobile component surfaces.\nThe electronic camera system is embodied as image-recognition system, which detects ahead of time the shape and position of automobile component surfaces to be measured optically with a spectrophotometer. The system carries out a form and position identification with the aid of the electronic databank, embodied as optical neuro-fuzzy structured image databank, and then initiates an optical-angle dependent spectrophotometer measurement, defined for this form and position, for the identified and classified automobile component surface, using predetermined measuring parameters such as wavelength, measuring angle, type of combination measuring—that is to say a measuring and scanning for color, gloss, layer thickness and wave.\nThe form of the automobile component surface is identified with the aid of electronic classes of automobile component surfaces which are stored object-related in the optical neuro-fuzzy structured image databank.\nThe neuro-fuzzy techniques are known in principle from the literature and have been used for years in different areas of the industry, that is for modeling, analysis, monitoring and control of industrial processes.\nIn contrast to the above methods, the method according to our invention distinguishes itself in that it is faster, meets more comprehensive object-specific measuring requirements and allows the classifying of measuring tasks with respect to the automobile component surfaces.\nThe present invention uses an optical neuro-fuzzy structured image databank in which the automobile component surface images are stored together with the associated measuring techniques. A camera image for comparing the image-patterns on the automobile component surface permits an allocation and/or classification of the measuring object and, following the object detection, controls the object-specific measuring technique of the angle-dependent spectrophotometer for the integral color, gloss, layer thickness and wave measuring with respect to the identified automobile component surface.\nDeviations in color, gloss, layer thickness and wave of the measured automobile component surfaces, including different automobile components, are stored object-specific and component-specific in a computer-aided optical quality databank, called a CAOQ databank, for the optical characterization. This CAOQ databank computes and administers as databank logical links, measured integrally for the various automobile component surfaces at different object points to obtain color, gloss, layer thickness and wave data and compares these data to required specified datasets for desired values predetermined by the automobile manufacturer for the automobile component surfaces as target values and tolerances. The data are transmitted with standardized statistic methods, outlier analyses, graphic representations of color, gloss, layer thicknesses and wave differences to a test station for surface measuring technology, via intelligent neuronal net, and can be visually displayed for an operator. Thus, the operator at the test station for surface measuring technology can reach a comprehensible decision that can be implemented with respect to the quality of individual, measured automobile component surfaces.\nAt least two mobile on-line and off-line monitoring methods of the above-described type are advantageously used at different locations during a production and are linked via intelligent neural net to the control station for surface measuring technology.\nIn a further step, the data and decisions recorded by the control station for surface measuring technology are transmitted to the production planning system, called PPS, via data transmittal and stored electronically for the specific automobile component. In particular, the production planning system records deviations in color, gloss, layer thickness and wave and the values determined with this method are then transmitted back electronically via the net to the production unit for automobile painting/coating, so that the detected deviations can enter into corresponding changes to the automobile painting/coating processes.\nIn an additional processing step, the transport system for the following automobile component surface to be measured is clock-pulse actuated.\nThe data are furthermore transmitted via network to a client-information system which transmits data on-line from the supplier to the automobile manufacturer.\nThe invention is explained in further detail with the exemplary embodiment shown in the drawing. This embodiment shows a schematic representation of a device for realizing the method for the mobile on-line and off-line monitoring of colored and high-gloss automobile component surfaces. Automobile component surfaces 1a-1c are produced in an automobile component production unit, not shown in further detail herein, and are then transported with the aid of a transport system 18 to the mobile on-line and off-line monitoring. The automobile component surface 1b, which is measured for example, is optically recorded with a camera 2. A pixel image signal for the automobile component surface 1b and its form is present at the output 3 of the camera 2. This information, which contains a pixel image that shows exact details of the surface image of the automobile component surface 1b, is transmitted to an optical neuro-fuzzy structured image databank 4 in which automobile component surface classes are electronically stored. In this optical neuro-fuzzy structured image databank 4, the real automobile component surface images are compared to the automobile image classes, stored therein, of the image databank 4. Following the identification of the automobile component surface 1b, the measuring technique associated with the object-class for the on-line and off-line monitoring of the colored and high-gloss automobile component surface 1b is initialized radio-wave supported 5a for the optical surface determination of color, gloss, layer thickness and wave and the angle-dependent spectrophotometer 6 is started with object-specific settings for the scanning with the measuring beam 7a. At the same time, the computer-aided optical quality databank 11 is initialized via radio wave 5b for the detector signal 10 data recording. The measuring beam 7a, formed with polarized light of different wavelengths from the angle-dependent spectrophotometer, then scans the automobile component surface 1b. The de-polarized measuring beam 7b, which is reflected by the automobile component surface 1b, travels to the angle-dependent spectrophotometer 6 with detector unit 8, a charge-coupled diode [CCD] array.\nThe output 9 of detector unit 8 carries a detector signal 10 for the surface conditions relating to color, gloss, layer thickness and wave. This detector signal, which contains reflection, interference, polarization and phase information and includes the values for color, gloss, layer thickness and wave for the automobile component surface 1b, is transmitted to a computer-aided optical quality databank 11, called CAOQ, which is initialized radio-wave supported 5b at point 4. This CAOQ databank 11 computes, compares and administers the detector signals 10 for different object points on the automobile component surface 1b and generates color, gloss, layer thickness and wave data that are compared to required data sets specified by the automobile manufacturer. These data are transmitted to different addressees via an intelligent neuronal network 12, which links at least two mobile on-line and off-line monitoring techniques 13 of the above-described type to different addressees and can be visualized on at least one control station for the surface measuring technology 14. The data recorded by the control station for surface measuring technology 14, are transmitted in the following step via electronic network to the production planning system 15, called PPS, and are then electronically stored for the specific automobile component. The deviations recorded in the production planning system 15 relating to color, gloss, layer thicknesses and wave data are transmitted via a different electronic network to the production unit for automobile painting/coating 16 and/or the client information system 17, so that necessary measures relating to automobile painting/coating can be initiated in unit 16. In a further step, the transport system 18 for the automobile components is controlled by the production planning system 15 and is clocked in time, so that the following automobile component surface 1c can be recorded and measured in accordance with the above-described method."}
{"text": "The operating speed of electronic devices such as a computer continues to increase.\nFor this reason, in recent years, as described, for example, in Japanese Laid-open Patent Publication No. 2000-332301, development of a photoelectrical composite substrate is being advanced for achieving communication between LSI or IC chips included in the electronic devices, by using optical signals. In a telecommunication field such as Internet, communication using optical signals is performed as described, for example, in Japanese Laid-open Patent Publication No. 2004-85756.\nWhen a photoelectric composite substrate is manufactured, for example, a drive element which drives an optical element and a package equipped with the optical element are prepared beforehand. The package is then mounted on a printed circuit board with the position of the optical element aligned with the position of an optical waveguide mounted on the printed circuit board. The package is mounted on the printed circuit board via an electrical connector such as a solder ball.\nWhen an electrical connector such as a solder ball has an electrical contact, an electrical signal may be transmitted therethrough. On the other hand, when an optical axis is not at a predetermined position, the connecting part between the optical element and the optical waveguide does not allow an optical signal to be transmitted. Thus, the positional precision demanded for the connecting part between the optical element and the optical waveguide is higher than the positional precision demanded for the electrical connector such as a solder ball. However, when the package equipped with the optical element is mounted on the printed circuit board, connection failure may occur between the optical element and the optical waveguide because of inevitable positional displacement therebetween caused by heat deformation or the like at the time of reflow.\nThus, it is an object of the present application to provide a method of manufacturing a photoelectric composite substrate, the method enabling relative positional precision between the optical element and the optical waveguide to be improved."}
{"text": "1. Field of the Invention\nThe present invention is in the field of electronic sensors and pertains particularly to an improved reflectance sensor usable in various commercial and consumer applications.\n2. Discussion of the State of the Art\nIn the field of sensory devices, more particularly electronic sensors, proximity sensing and motion detection are regimens that provide contact-less control and object detection useful in a large variety of consumer, industrial, and security applications.\nDevelopment of various electronic technologies for proximity sensing has occurred and development continues. Each accepted technology has provided one or more advantages depending on the specific application of those technologies. These techniques can be classified in terms of the operating principle of the device versus the detection medium used, whether light, radio waves, or the like.\nTo provide one example, a well-known proximity sensor based on measuring an echo transit time that uses radio waves as a medium is called a radar. Another echo-principled proximity sensor that uses sound as a medium is called a sonar sensor. Still, another echo-principled proximity sensor that uses light as a medium is called a light-imaging detection and ranging (LIDAR) sensor. Although these classic echo-based sensors provide relatively accurate distance and speed information, they can, depending on the application, be expensive, bulky, may consume high power, and/or may require very high frequency technologies to be successful.\nIn contrast to the sensor device types described above, there are less expensive and much smaller proximity sensors that use field disturbance techniques to detect proximity. These sensors may be classed as either passive or active sensors. These types of sensors detect proximity base on changes in a field caused by interactions with a detected object.\nOf the above-described sensors, passive field-disturbance detectors use background radiation or the emissions of an object as the source of the field. To exemplify, a sound-activated switch may be provided in a “smart toy” to sense when a human is talking nearby. In another case of what would be termed a passive sensor, an infrared motion detector may sense changes in the infrared background due to movement of objects irradiating infrared above the background because of their higher temperature. Such passive infrared sensors are widely used in high-volume applications for alarm systems, automatic light turn-on sensors, and the like.\nAn active proximity sensor may detect changes in field disturbances caused by an active source. One example of such an active sensor may be a magnetic field-disturbance proximity sensor (metal detector) that senses changes in a local alternating current (AC) magnetic field due to eddy current or magnetic characteristics of metal objects. Another well-known example is that of a “stud finder” used in carpentry. A stud finder is an AC electrostatic field disturbance sensor that senses the change in dielectric constants between air and wood as the finder is slid across overlying plaster.\nStill another example of an active field disturbance detector is a proximity sensor that detects changes in optical reflectance, avoiding technical difficulties of nanosecond time resolution necessary for RF or optical echo transit time devices. Optical reflectance proximity sensors have certain key advantages over other types of proximity sensors for sensing objects in the range of 1 to 100 centimeters (cm). The advantage may be due to the fact that they typically use a small LED as a light source and a small photodiode driving a receiver circuit. Another advantage is that they can be small enough to fit in miniature electronic devices. Even these types of active reflectance sensors, although fairly small and robust, have been too expensive for consumer applications and consequently have found applications mostly in the industrial and commercial markets.\nPassive infrared sensors are the most common in consumer application due to lower cost. However, they are too large for many consumer products because they require a collecting lens and a photodiode of several centimeters in diameter to gather a sufficient signal to sense changes in the ambient background visible to the sensor. Moreover, a major handicap of such sensors is a lack of reliability earmarked by spurious triggering and by failure to trigger. For example, large, far-away objects or sudden temperature changes may trigger them and they can completely miss radiation-neutral objects.\nReflectance sensors have been developed that overcome the handicaps described immediately above even though they do not inherently measure distance (as echo transit-time devices do). Reflectance sensors can detect an object moving into certain ranges unambiguously because of a strong fourth-power decrease of reflected light from the object sensed. A 20% change in distance will cause approximately a 100% change in reflected signal. For objects with a 10-to-1 variation in reflectance between them, this amounts to less than 50% difference in detection range.\nSmall, short range (1 cm to 2 m) reflectance proximity sensors are more useful than passive infrared sensors in many applications. However, widespread use in consumer applications is not apparent because of higher cost factors.\nIn general, proximity-sensor applications break down into two broad functional groups: (1) those that provide an on-off function and (2) those that provide analog or digital proportional information. Examples in the first category include: automatic flushing for a public lavatory; an automatic doorbell that detects a person passing through a door; or an object sensor on an automatic production line. Examples in the second category are: an automobile bumper warning indicator that puts out an audible warning signal whose pitch is proportional to the distance from an obstacle; a toy car that slows down when it approaches an object and steers away from it; or a light switch that can be activated and dimmed by waving one's hand near it.\nProximity sensors that provide on-off functions can often serve as replacement for switches that are either operated manually or by some other machine function. In both cases, the electronic proximity switch will be more reliable than a standard mechanical switch especially if a very high number of cycles occur over the life of the switch. But since most proximity sensors require significant amounts of power, they are not normally used for power switches on battery-powered products. Of course, like normally open or normally closed switches, a proximity sensor can provide either an “on” or “off” function when an object moves in or out of proximity.\nFurther to the above, some proximity-sensor applications also need to measure the ambient background light such as a proximity-activated security light. Proximity sensors that provide analog proportional information can functionally replace analog controls allowing smarter processing of proximity information for more complex applications. However, unlike on-off proximity sensors, proportional sensors generally need to interface to a microprocessor.\nThe inventors are aware of a method taught by Holcombe (U.S. Pat. No. 5,864,591) for using circuitry to reduce feedback in an infrared data receiver. The method includes a circuit that is configured as an infrared receiver including an automatic gain control (AGC) circuit where the AGC is isolated from the input to the receiver in response to the output signal from the receiver in order to suppress the effect of feedback from the output signal to the input of the receiver.\nThe inventors are also aware of an enhancement to the method taught by Holcombe (U.S. Pat. No. 6,240,283). The enhancement includes a method and apparatus for controlling the input gain of a receiver whereby the input gain is controlled by sampling an amplified data signal during a time interval when a positive-going feedback transient from an output terminal of the receiver to an input terminal of the receiver is not present in the amplified data signal.\nIn this enhanced circuit, an input amplifier has variable gain determined by a gain control signal, a comparator which compares the amplified data signal from the input amplifier to a detection threshold voltage to produce a demodulated data signal and an analog delay circuit which delays the amplified data signal by a predetermined time interval to produce a delayed data signal. The method is enabled by a switch that is driven by the demodulated data signal to sample the delayed data signal for input to an automatic gain control circuit. The automatic gain control circuit compares the sampled delayed data signal to an automatic gain control threshold potential and rectifies and integrates the resulting waveform to produce the gain control signal.\nIn one application, the data signal is amplified by a gain factor and the signal is then compared to a detection threshold voltage to produce a demodulated data signal. The amplified signal is also delayed to produce a delayed data signal. The delayed signal is sampled using the demodulated data signal to produce a sampled data signal that is used to adjust the gain factor in the amplifier. Although the technique described by Holcombe may provide some feedback immunity from a detected infrared receiver output to a photodiode in a IRDA communications receiver, there may also be significant feedback from the LED driver to the photodiode. For example, the LED driver produces both a voltage and an inductive current transient when it initially turns on. The voltage transient can couple to the photodiode input via wire-bond and PCB-trace capacitances. The inductive transients can couple to the photodiode input through ground and power-supply traces. These transients (voltage and inductive) may produce spurious signals that mask low-level reflectance signals, thus limiting the minimum detectable signal levels in the communications receiver circuit.\nIt has occurred to the inventor that with some innovative enhancement to the feedback immunity techniques described in U.S. Pat. No. 5,864,591, and in U.S. Pat. No. 6,240,283, a low cost optical reflectance proximity sensor could be provided that could overcome the problems associated with the relevant art described above.\nTherefore, what is clearly needed in the art is a reliable and sensitive optical reflectance proximity sensor that is very small and inexpensive to manufacture. Such a sensor would consume very little power, would not be required in all applications to interface with a microprocessor, and could be implemented in some analog output applications without the complexity and cost of using a standard digital-to-analog (DAC) converter."}
{"text": "Interventional cardiologists incorporate a variety of diagnostic tools during catheterization procedures in order to plan, guide, and assess therapies. Fluoroscopy is generally used to perform angiographic imaging of blood vessels. In turn, such blood vessel imaging is used by physicians to diagnose, locate and treat blood vessel disease during interventions such as bypass surgery or stent placement. Intravascular imaging technologies such as optical coherence tomography (OCT) are also valuable tools that can be used in lieu of or in combination with fluoroscopy to obtain high-resolution data regarding the condition of the blood vessels for a given subject.\nIntravascular optical coherence tomography is a catheter-based imaging modality that uses light to peer into coronary artery walls and generate images for study. Utilizing coherent light, interferometry, and micro-optics, OCT can provide video-rate in-vivo tomography within a diseased vessel with micrometer level resolution. Viewing subsurface structures with high resolution using fiber-optic probes makes OCT especially useful for minimally invasive imaging of internal tissues and organs, as well as implanted medical devices such as stents.\nStents are a common intervention for treating vascular stenoses. It is critical for a clinician to develop a personalized stent plan that is customized to the patient's vascular anatomy to ensure optimal outcomes in intravascular procedures. Stent planning encompasses selecting the length, diameter, and landing zone for the stent with an intention to restore normal blood flow to the downstream tissues. However, flow-limiting stenoses are often present in the vicinity of vascular side branches. Side branches can be partially occluded or “jailed” during deployment of a stent intended to address a stenosis in the main vessel. Since side branches are vital for carrying blood to downstream tissues, jailing can have an undesired ischemic impact and also can lead to thrombosis. The ischemic effects of jailing are compounded when multiple side branches are impacted or when the occluded surface area of a single branch is increased.\nMetal stent detection methods typically detect individual stent struts by detecting shadows cast by the struts onto the blood vessel wall, followed by detecting the location of the struts within the detected shadows. However, struts over jailed side branches are difficult to detect via this method. Side branches appear as large shadows in images because the scan line can be perpendicular to the side branch opening. As a result, it is difficult or impossible to detect strut shadows overlying side branches. Consequently, jailing struts are easily missed by the shadow based detection methods.\nThe present disclosure addresses the need for enhanced detection of jailing stent struts."}
{"text": "The present invention relates generally to the fabrication of thin-walled articles of tungsten-nickel-iron alloy and more particularly to a method of fabricating such articles wherein the articles are of near theoretical density and possess essentially uniform properties.\nTungsten has proven to be a particularly useful material for various industrial applications and in the construction of nuclear reactors. The tungsten metal can be placed in a somewhat ductile form without seriously detracting from the desirable properties by alloying it with nickel and iron. An alloy composition found to be particularly useful is formed of essentially 95 weight percent tungsten, 3.5 weight percent nickel and 1.5 weight percent iron. However, alloy compositions with a nickel concentration in a range of about 2.1 to 7.0 weight percent and an iron concentration in a range of about 0.9 to 3.0 weight percent iron have proven to be useful in essentially the same areas as the 3.5 weight percent nickel and 1.5 weight percent iron compositions. Structures of the ductile tungsten-nickel-iron alloys are generally prepared by employing conventional powder metallurgical procedures. These procedures usually comprise the blending desired weight percents of elemental tungsten, nickel, and iron alloy powders, pressing the blended powders into a compact of the desired configuration and thereafter sintering the compact at a sufficiently high temperature to convert the minor phase of the alloy into a liquid to provide structural integrity to the sintered structure. Final densification, microstructure and properties of the structure are dependent upon the pressing and liquid phase sintering operation.\nWhile tungsten-nickel-iron alloy structures of various configurations have been successfully fabricated by cold pressing and sintering the resulting compact, there have been some difficulties associated with the fabrication of structures having relatively thin walls, such as open cylinders and the like, because of the required liquid phase sintering operation. More specifically, the pressed compacts shrink approximately 40 volume percent during the liquid phase sintering step so that the thin-walled articles are usually subjected to extensive distortion and stresses which lead to deleterious cracking and other stress associated problems. Some techniques have been employed to overcome the distortion and cracking problems such as the use of multiple sintering steps and special sintering mandrels. However, even with such techniques, uniform, physical and mechanical properties are seldom achieved. For example, the density and tensile strength of the sintered structure usually vary by greater than 1 percent and 4 percent, respectively, over the structure."}
{"text": "1. Field of the Invention\nThe present invention relates to a photovoltaic module, and particularly relates to a thin film solar cell module of see-through type.\n2. Description of Related Art\nSolar energy is a renewable energy, which causes no pollution. It has been the focus in the development of environmental-friendly energy as an attempt to counter the problems such as pollution and shortage of fossil fuels. Herein, solar cells can be used to directly convert solar energy into electrical energy, which becomes a very important research topic now.\nCurrently, mono-silicon and poly-silicon cells account for more than 90% of the solar cell market. However, manufacturing these types of solar cells requires silicon chips with thickness of 150˜350 micrometers, which increases the production costs. Furthermore, the raw material of solar cells is high-purity silicon ingot. Due to the significant increase in the consumption of silicon ingot, it is being depleted by day. Hence, thin film solar cells have become the new direction in the research and development of solar energy. Thin film solar cells are suitable for mass production and have the advantages of lower production costs and simpler module fabricating process.\nFIG. 1 schematically illustrates a conventional thin film solar cell module. As shown in FIG. 1, a thin film solar cell module 150 comprises a glass substrate 152, a transparent electrode 154, a photoelectric conversion layer 156, and a metal electrode 158. Herein, the transparent electrode 154 is disposed on the glass substrate 152. The photoelectric conversion layer 156 is disposed on the transparent electrode 154 by position displacement. In addition, the metal electrode 158 is disposed on the photoelectric conversion layer 156 by position displacement and is in contact with the transparent electrode 154 underneath. In the thin film solar cell module 150, the photoelectric conversion layer 156 usually includes a p-i-n structure composed of a p-type semiconductor, an intrinsic semiconductor, and an n-type semiconductor. Light is transmitted through the bottom of the glass substrate 152 and is absorbed by the photoelectric conversion layer 156 to generate electron-hole pairs. The electron-hole pairs are then separated by an electric field established across the device to form a voltage and an electric current, which are transmitted by a conductive wire for loading. To enhance the efficiency of the cells, in the conventional thin film solar cell module 150, pyramid-like structures or textured structures (not shown) are formed on the surface of the transparent electrode 154 to reduce the reflection of light. The photoelectric conversion layer 156 is usually formed by using an amorphous silicon thin film. However, the band gap of the amorphous silicon thin film is usually between 1.7 eV and 1.8 eV, which merely absorbs sunlight of wavelength less than 800 nm. To increase the utility of light, usually a layer of micro-crystalline (or nano-crystalline) thin film is stacked on the amorphous silicon thin film to form a p-i-n/p-i-n tandem solar cell. The band gap of micro-crystalline (or nano-crystalline) is usually between 1.1 eV and 1.2 eV, which absorbs sunlight of wavelength less than 1100 nm.\nIn the early times, it was costly and difficult to manufacture solar cells, and solar cells were only used in special fields such as astronautics. Now solar cells, which feature converting solar energy into electric energy, have become more widely used and applied. The applications of solar cells range from the use in apartments and high-rise buildings to that in camper vans and portable refrigerators. However, silicon wafer solar cells are not suitable for certain applications such as transparent glass curtains and buildings integrated with photovoltaic (BIPV). Thin film solar cells of see-through type are used in the aforesaid applications because they are energy-efficient and pleasing to the eye. Further, they accommodate more readily with our living demands.\nCurrently, techniques related to the thin film solar cells of see-through type and the methods for fabricating the same have been disclosed in some U.S. patents.\nU.S. Pat. Nos. 6,858,461 (6,858,461 B2) provides a partially transparent photovoltaic module. As shown in FIG. 2, a photovoltaic module 110 includes a transparent substrate 114, a transparent conductive layer 118, a metal electrode 122, and a photoelectric conversion layer disposed between the transparent conductive layer 118 and the metal electrode 122. Similarly, light is transmitted through the bottom of the transparent substrate 114. In the photovoltaic module 110, a laser scribing process is performed to remove a portion of the metal electrode 122 and a portion of the photoelectric conversion layer to form at least one groove 140 so as to achieve partial transparency of the photovoltaic module 110. However, the laser scribing process is performed at a high temperature. Due to such a high temperature, the metal electrode 122 easily forms metal residues or melts down and accumulates in the grooves, resulting in short circuits of the top and bottom electrodes. On the other hand, an amorphous silicon photoelectric conversion layer recrystallizes at such a high temperature and forms low resistant micro-crystalline (or nano-crystalline) silicon on the sidewalls of the grooves. Consequently, current leakage is increased, and the process yield and the efficiency of the solar cells are affected. In addition, pyramid-like structures or textured structures are usually formed on the surface of the transparent conductive layer 118 to enhance the efficiency of the cells. However, light transmittance is not effectively enhanced because the light transmitted through the bottom of the transparent substrate 114 would be scattered.\nIn view of the above, to achieve a certain level of light transmittance, larger portions of the metal electrode and the photoelectric conversion layer in a solar cell need to be removed. Please refer to Table 1, which lists the technical specifications of various thin film cells of see-through type manufactured by MakMax Taiyo Kogyo (Japan). According to Table 1, to increase light transmittance, larger portions of the metal electrode and the photoelectric conversion layer need to be removed to decrease the maximum output, efficiency, and fill factor (FF).\nTABLE 1TypeKN-38KN-45KN-60Size (mm)980 × 950980 × 950980 × 950Transmittance Rate (%)105<1Maximum Power Output (W)38.045.058.0Vpm (V)58.664.468.0Ipm (A)0.6480.6990.853Voc (V)91.891.891.8Isc (A)0.9721.0901.140Efficiency4.14.86.2FF0.430.450.55\nIn addition, a photovoltaic device is disclosed in U.S. Pat. Nos. 4,795,500 (4,795,500). As shown in FIG. 3, a photovoltaic device includes a transparent substrate 1, a transparent conductive layer 3, a photoelectric conversion layer 4, a metal electrode 5, and a photoresist layer 8. In the photovoltaic device, holes 6 are formed in the metal electrode 5, the photoelectric conversion layer 4, and even in the transparent conductive layer 3 to achieve transparency. Nevertheless, this patent utilizes a lithographic process which requires expensive facilities and increases the production costs. Additionally, if this patent utilizes a laser scribing process to directly form the holes 6, the problems of metal residue contamination and short circuit will occur to affect the process yield."}
{"text": "The present invention relates to a condenser for condensing exhaust steam of steam turbines of thermal power plants and atomic power plants.\nGenerally in a steam turbine plant, steam which has worked in a steam turbine and inflated is condensed by a surface contact type-condenser for recovery.\nFIG. 11 is a sectional view of one example of the above-described condenser. In a condenser body 1 into which exhaust discharged from the steam turbine not shown there are disposed cooling pipe bundles 2 of a number of cooling pipes which are extended in a first direction (which is perpendicular to the sheet of FIG. 11) and in parallelism with each other, whereby the exhaust from the steam turbine is heat-exchanged with cooling water, such as sea water, river water or others, on the surfaces of the respective cooling pipes and condensed to be drained.\nEach cooling pipe bundle 2 is divided in a plurality of pipe groups 2a, 2b, 2c, 2d, 2e, and the pipe groups are defined by partition plates 3, 4 so that the heat conducting pipe groups do not affect by the drain the heat-exchange of the other heat conducting pipe groups. Air cooling pipe groups for condensing residual energy of the steam 6, 7 are provided below the cooling pipe bundle 2. Partition plates 5a, 5b are provided between the cooling pipe bundle 2 and the pipe groups 6, 7 respectively. Gas discharging devices 8, 9 are provided respectively on the sides of the air cooling pipe groups 6, 7. Enclosure plates 10, 11 are provided respectively below the air cooling pipe groups 6, 7, and a sprinkler box 12 having a U-seal is provided between the enclosure plates 10, 11 and therebelow.\nDrain thus heat-exchanged and condensed by the cooling pipe bundle 2 as the steam flows is collected at the center of the pipe bundle 2 and flows into the below sprinkler box 12 through between the air cooling pipe groups 6, 7 enclosed by the enclosing plates 10, 11, then falls through the U-seal into a hot well 13 which is a lowermost part of the condenser body 1 and discharged outside through a drain exit 14.\nOn the other hand, uncondensed gases, such as steam which could not be condensed by the cooling pipe bundle 2, air, etc. flows through the air cooling pipe groups 6, 7 horizontally toward the outside of the condenser body to be discharged to the outside of the condenser body 1 through an air exhaust pipe 15 via the gas discharging devices 8, 9.\nBecause the air exhaust pipe 15 is connected to the outside of the air cooling pipe groups 6, 7 as described above, drain generated in the air cooling pipe groups 6, 7 tends to intrude the air exhaust pipe 15. For the purpose of prohibiting the drain intruding the air exhaust pipe 15 as described above from residing in the air exhaust pipe to return to the sprinkler box 12, the air exhaust pipe includes a vertical piping or an inclined piping portion to lead the uncondensed gas upward, whereby exhaust of the uncondensed gas is smoothed, and downstream machines and instruments are protected from erosion and corrosion. Accordingly it is necessary to arrange the air exhaust pipe extended upward along the sides of the cooling pipe bundle 2.\nThe drain once intruded the air exhaust pipe is sometimes carried against the gravity by the uncondensed gas in the air exhaust pipe when the uncondensed gas has a high flow rate. In addition, pressure loss increase much affects achievement of the condenser. Accordingly, thick pipes are used for low flow rates.\nHowever, the sides of the cooling pipe bundle 2 are places where the steam which has flowed from upward flows into the cooling pipe bundle 2. In these places the flow passage has a most restricted area by the cooling pipe bundles 2, and the steam has a highest flow rate.\nAccordingly it causes pressure loss increase to dispose the air exhaust pipe 15, which is to be a barrier, in such a high flow rate area. As a result of that, the outlet pressure increase of the turbine occurs. Resultant problems are that effective use of thermal energy is affected, and others. The air exhaust pipes 15 must be arranged extended between the cooling pipe bundles and between the cooling pipe bundle and the condenser body, which results in the increased width of the bottom of the condenser body 1. Problems are that the condenser cannot be compact, and others.\nIn view of the above-described problems, the present invention was made, and an object of the present invention is to provide a condenser which can prevent pressure loss of steam flow due to the air exhaust pipe, prevent thermal energy loss, and can be compact."}
{"text": "1. Field of the Invention\nThe present invention relates to a display apparatus which provides a capability of displaying a color image.\n2. Description of the Related Art\nAs a color image display apparatus, a liquid crystal display has been proposed which includes a liquid crystal panel served to control a light transmittance of each pixel and a backlight unit combined therewith so that a color image may be displayed.\nTo display a color image, it is necessary to include at least three primary color components of RGB (Red, Green and Blue) in the backlight and locate sub-pixels each of which has one of at least RGB color filters as the pixels composing the liquid crystal panel. This arrangement makes it possible to control a light quantity on the overall range of a wavelength. Herein, the sub-pixel provides any one of RGB color filters and corresponds to a minimum unit of transmittance control. The pixel, termed herein, designates a combination of three sub-pixels of the RGB. Lots of pixels are ranged on the screen plane of the display apparatus. The other display apparatuses of a CRT (Cathode-ray Tube) system, a plasma system, a projector system and so forth have the same fundamental principle of display as that of the liquid crystal display apparatus, that is, those apparatuses also display a color image by properly arranging the pixels.\nIn the meantime, illumination brightness is considered as an environmental condition of a place where the display apparatus is located. An observer of a display screen watches an image in the synthesized light of ambient light reflected on the display screen and the light displayed by the display apparatus itself.\nLetting a contrast R be a ratio of a maximum to a minimum of brightness of a display screen, the relation between the display light and the reflected light is represented as follows:R=(Maximum Display Light Quantity+Reflected Light Quantity)/(Minimum Display Light Quantity+Reflected Light Quantity)\nIn general, as the contrast R becomes greater, a visibility is made better. Herein, the maximum display light quantity designates a display light quantity corresponding with a maximum value of a display signal and the minimum display light quantity designates a display light quantity corresponding with a minimum value of a display signal. To improve the contrast, it is effective to make the maximum display light quantity greater or the reflected light quantity smaller.\nTo manage the contrast of the display screen, a method has been known in which a light sensor for sensing a brightness of ambient light is prepared so that an intensity (luminance) of the display light may be variably set according to the output of the light sensor.\nThe JP-A-2006-106294 discloses a technique of varying a luminous quantity of the backlight unit depending upon the ambient light, for example, in a gloomy room or a bright outdoor place. Concretely, during daylight, the technique is caused to increase the maximum display light quantity of the display screen by raising the luminous quantity of the backlight unit based on the output signal of the light sensor. This technique makes the contrast R higher and thereby the visibility better.\nThe US 2005/0225562 proposes a technique of adding a sub-pixel of W (White) to a combination of three RGB sub-pixels of RGB in order to improve the luminance of the display panel itself. The W sub-pixel provides no color filter, so that it has a high light transmittance and thus is effective in improving the luminance. Concretely, when the conventional pixels each composed of the RGB sub-pixels are compared with the pixels each composed of the RGBW sub-pixels on the same area, the area ratio of the sub-pixels of the former to the latter is made to be 4:3. The RGB color filters cut a wavelength distribution of a light source into one third, while the W color filter transmits a light quantity of a light source as it is. In light of this relation, the ratio of the maximum display light quantity of the RGB panel to the RGBW panel is made to be ((4+4+4)/3):((3+3;3)/3+3×1)=1:1.5\nAs a method of generating a RGBW signal to be used for driving the RGBW pixels, the following process has been proposed. The minimum values of the input color signals of the RGB are set to a W=MIN (R, G, B) signal or the values derived by subtracting W from the RGB color signals are newly set to a R′ G′ B′(R′=R−W, G′=G−W, B′=B−W) signal. By multiplying a proper amplification factor by the W signal, it is effective to improve the luminance.\nFurther, the technical background of the color reproducibility is discussed in detail in “Chromatic Science Handbook, Second edition, edited by Chromatic Academic Society of Japan, published by Tokyo University Publication, 1998”."}
{"text": "1. Field of the Invention\nThe invention relates to a droop circuit, and more particularly to a multi-phase DC-DC converter with a droop circuit.\n2. Description of the Related Art\nIn general, a droop circuit can control a swing of an output voltage in a DC-DC converter. FIG. 1A shows an output voltage of a DC-DC converter without a droop circuit, and FIG. 1B shows an output voltage of a DC-DC converter with a droop circuit. As shown in FIG. 1A, the output voltage is operated at a normal voltage Vnom except for a time t1 and a time t2. Transients of the output voltage are generated due to an output current I1 which is varied with a load of the DC-DC converter, and rapidly increases and decreases at the time t1 and the time t2. Referring to FIG. 1B, an output current I2 of the DC-DC converter with a droop circuit rapidly increases and decreases at a time t3 and a time t4. The output voltage is operated at a minimum voltage Vmin during a duration between the time t3 and t4. However, the output voltage is operated at the normal voltage Vnom outside of the duration between the time t3 and t4. Thus, for a DC-DC converter, transients of the output voltage are avoided by the droop circuit.\nFIG. 2 shows a conventional multi-phase switching regulator disclosed in U.S. Pat. No. 6,683,441. In FIG. 2, an amplifier circuit 28 generates a voltage Vcs according to an output voltage Vout and a summing voltage of a summing node 26. Then, the output voltage Vout is subtracted from the voltage Vcs to generate a droop voltage Vdroop by a summation circuit 30. Thus, the multi-phase switching regulator needs the summation circuit 30 to obtain the droop voltage Vdroop, and the output voltage Vcs of the amplifier circuit 28 is equal to the output voltage Vout plus the droop voltage Vdroop. FIG. 3 shows a conventional droop amplifier circuit disclosed in U.S. Pat. No. 7,064,528 for generating a droop voltage VDROOP. In FIG. 3, a positive polarity (+) of the droop voltage VDROOP is provided by an output of an amplifier A2, and a negative polarity (−) of the droop voltage VDROOP is provided by an output node of a multi-phase DC-DC regulator. Hence, an output voltage of the amplifier A2 is equal to an output voltage VOUT of the multi-phase DC-DC regulator plus the droop voltage VDROOP."}
{"text": "1. Field of the Invention\nThe instant disclosure relates to an auxiliary detachable assembly for electronic devices; in particular, to a detachable assembly and the memory module using the same which can prevent electrostatic interference.\n2. Description of Related Art\nAs the volume of electronic devices continues to miniaturize, the processing speed and operational functions are gradually becoming faster and more powerful, respectively. As a result, mainboards and electronic components in electronic devices are not only required to be miniaturized, but are also required to provide even higher operational speed in order to meet users' demand. As the data processing power increases, the demand for higher capacity and speed from memory modules of computers increases.\nDynamic random access memory (DRAM) is presented in the form of a module within a computer. In other words, DRAM module includes integrated circuits and printed circuit boards. Users usually first align the printed circuit board of the DRAM module against the memory slot, then apply a force on the top portions of the DRAM module such that the printed circuit board is inserted into the memory slot, thus, the computer system can operate normally.\nHowever, static electricity tends to store in the hands of users and is transmitted to the DRAM module which can lead to abnormal operations thereof. In severe cases, electronic components may malfunction or even be damaged. Moreover, the surface area on top of the DRAM module in which users can apply force upon is relatively small. Specifically, the force applicable area is substantially the thickness of the printed circuit board. As a result, users are more likely to misalign the DRAM module with the memory slot, thus, affecting normal operations, or even worse, damaging the DRAM module.\nTo address the above issues, the inventor strives via associated experience and research to present the instant disclosure, which can effectively improve the limitation described above."}
{"text": "In Japan Patent Application Nos. 67,013/84 and 69,929/84, the present inventor proposed vibratory apparatuses which are directed for hypnosis and dehypnotization, as well as for therapy of myalgia and stiffness in the shoulder.\nThese apparatuses, however, have the disadvantage that they consume a relatively large amount of electricity to radiate heat when used for a long duration, because these apparatuses are of forced vibration-type wherein a magnetic coil and an iron core are attached to the same iron plate."}
{"text": "Plasma arc torches are widely used in the cutting, and marking of materials. A plasma torch generally includes an electrode and a nozzle having a central exit orifice mounted within a torch body, electrical connections, passages for cooling, and passages for arc control fluids (e.g., plasma gas). Optionally, a swirl ring is employed to control fluid flow patterns in the plasma chamber formed between the electrode and nozzle. In some torches, a retaining cap can be used to maintain the nozzle and/or swirl ring in the plasma arc torch. The torch produces a plasma arc, a constricted ionized jet of a gas with high temperature and high momentum. Gases used in the torch can be non-reactive (e.g., argon or nitrogen) or reactive (e.g., oxygen or air). In operation, a pilot arc is first generated between the electrode (cathode) and the nozzle (anode). Generation of the pilot arc can be by means of a high frequency, high voltage signal coupled to a DC power supply and the torch or by means of any of a variety of contact starting methods.\nOne category of hand held plasma arc torch systems include a manual gas control knob on the control panel of the power supply or power supply housing. Before cutting a workpiece, an operator is required to manually adjust the gas pressure or gas flow rate based on the process parameters set forth in a cut chart. The operator manually adjusts the gas pressure or flow rate for each type of cut and therefore, constantly refers to the cut chart for the appropriate gas pressure or flow rate. Moreover, if the operator inadvertently inputs an incorrect gas pressure or flow rate, the plasma arc torch can operate incorrectly or can operate inefficiently.\nAnother category of hand held systems eliminate the gas control by automatically setting the gas pressure based on the user selected current level and mode (i.e., gouging or cutting). This category of hand held plasma arc torches does not provide the operator with any flexibility in setting the gas pressure beyond the preset automated systems. Therefore, if the operator determines that the gas pressure or flow rate should be changed due to a changed operating parameter or to optimize the plasma arc torch, the operator does not have the flexibility to make these operational and/or optimizing adjustments."}
{"text": "This invention relates to cleaning systems and apparatus, in particular a vehicle used for pumping waste from a pond, pit or other area.\nIn many industrial settings certain areas have been used to dump various solid and liquid wastes. These areas are often ponds or pits and contain a mixture of water, nonaqueous liquids, suspended solids, flowable wastes and solid trash such as bottles, old tires and used drums.\nIt has become painfully evident that such waste sites can create a severe environmental problem due to leaching of substances into the earth. Even if leaching is not a problem, such sites can create serious problems if simply filled in with soil once the site is no longer needed. Accordingly, there has been an increased interest in cleaning these sites, collectively called sludge ponds in this application. However, because of the physical makeup of sludge ponds, substantial obstacles hinder cleanup. There are problems associated with the process of physically removing and separating the water, muck, trash and other material, collectively termed sludge, found in the sludge pond. The wide variety of toxic and noxious materials often dumped into or which leak into sludge ponds creates substantial health risks to those removing the sludge."}
{"text": "There is a constant rise in demand for artificial meniscal grafts mimicking native articular tissue to be used for surgical treatment of meniscal lesions. In Europe alone over 400,000 surgical cases involving the meniscus are being performed annually, and over 1 million similar cases are treated in the United States. By far meniscectomy is known to be the most common surgical procedure performed in the orthopedic field today. The current therapeutic strategy for this type of meniscus tears is either partial or subtotal meniscectomy, with only a small percentage being successfully repaired but finally leading to osteoarthritis of the knee with time (Fairbank, 1948; Englund et al., 2003).\nA functional intact meniscus is of paramount importance for homeostasis of the knee joint. It helps perform complex knee joint biomechanics, in load bearing, load transmission, shock absorption, joint stability and joint lubrication. However, due to lack of vasculature, human meniscus has a poor healing potential. Blood vessels are reported to be present only in the outer 10-30% of the meniscal body and can be sutured successfully with a high success rate (Englund et al., 2003; Buma et al., 2004). In contrast, majority of these meniscal tears are situated in the inner avascular zone lacking spontaneous healing process and hence be resected (Kohn et al., 1999). Removal and/or damage of all this important anatomical structure eventually leads to degenerative changes of the articular cartilage, osteoarthritis and subsequent clinical symptoms due to increased peak stresses (Fairbank, 1948; Cole et al., 2003; Chatain et al., 2003; Englund et al., 2003). It has been estimated that cartilage volume loss after meniscectomy is at 4% per year and is known to be more pronounced in the lateral compartment as compared to medial compartment (Verdonk and Kohn, 1999).\nTo this problem, meniscus allo/autograft transplantation represents a potential tissue engineering solution for the symptomatic, meniscus deficient patient to substitute for lost meniscal tissue to prevent cartilage degeneration, relieve pain and to improve function. The strategies included delivery of potent cells to the defect site for repair including chondrocytes, fibrochondrocytes and stem cells (Peretti et al., 2004; Izuta et al., 2005; Port et al., 1996). The other strategy being direct replacement of defective tissue in part or as a whole has also been carried out using both natural and synthetic scaffolds, including collagen-based grafts, subintestinal submucusa, cell free hydrogels, degradable porous foams, macro- and microporous polymeric meshes to improve immediate or long term outcomes (Buma et al., 2004; Stone et al., 1992; Cook et al., 2006 a; Setton et al., 1999; Sweigart et al., 2001; Kobayashi et al., 2005; Kelly et al., 2007; Van Tienen et al., 2002; Heijkants et al., 2004; Cook et al., 2006 (a, b). In the past, a variety of these materials have already been reported for cartilage tissue engineering including, poly-glycolic acid (PGA), poly-L-lactic acid (PLA), copolymer poly-lactic-co-glycolic acid (PLGA) and alginate (Grande et al., 1997; Freed et al., 1993 a,b; Paige et al., 1996; Marijnissen et al., 2002; Ma et al., 2003). However, these materials have intrinsic limitations, including inflammation in vivo in the case of the polyesters and rapid degradation and high swelling in the case of collagen, which can limit their use (Cancedda et al., 2003; Athanasiou et al., 1996; Wakitani et al., 1994; Meinel et al., 2004 a,b). In terms of meniscus shape, a PGA spun matrix was used in a rabbit model but failed to recapitulate the complex internal meniscus architecture (Kang et al., 2006). Additional efforts have focused on mimicking the native mesh-like meniscus architecture using cell alignment on biodegradable electrospun fibers for enhanced biomechanics (Baker and Mauck; 2007; Baker et al., 2009). Many of the above studies employed in vivo animal models to show chondroprotection by the implant, but with a low success rate due to failure to mimic the complex internal architecture and biomechanics of the native meniscus.\nIn order to develop a functional tissue engineered meniscus, mimicking its complex internal architecture is most important. In this regard, none of the approaches previously reported have successfully recapitulated the complex native meniscal multiporous and aligned structure as a single meniscus wedge shaped unit to completely and/or partially eliminate cartilage degeneration. Thus, in order to mimic the meniscus in a tissue engineered approach, understanding its structural and functional components is important. Menisci are wedge-shaped semi-lunar discs present in duplicate in each knee joint which are attached to the transverse ligaments, the joint capsule, the medial collateral ligament (medially) and the menisco-femoral ligament (laterally) (McDevitt and Webber, 1990; Sweigart and Athanasiou, 2001). An extensive scanning electron micrograph study of the human meniscus by Peterson and Tillmann showed 3 distinct zones comprising of outer finer meshwork, middle broader mesh like fibrous structure and bottom most aligned collagen bundles in laminar orientation (Petersen and Tillmann, 1998). This particular aligned laminar orientation of fibers along with mesh structure within was reported to contribute maximally for its high intrinsic tensile and compressive properties of native meniscus (Sweigart and Athanasiou, 2001; Tissakht and Ahmed, 1995; Petersen and Tillmann, 1998). As a fibrocartilaginous structure, the meniscus has characteristic of both fibrous (outer region) and cartilaginous (inner region) properties (O'Connor, 1976; Petersen and Tillmann, 1998). Knee meniscal fibrocartilaginous tissue contains mainly water (72%), collagens (22%) and glycosaminoglycans (0.8%) (Proctor et al, 1989; Herwig et al, 1984). Of the total collagen content, Type I collagen accounts for over 90%. The remaining 10% are meniscal collagens Type II, III and V collagen (Eyre and Wu, 1983; McDevitt and Webber, 1990). It has been shown that peripheral two-thirds of the meniscus solely consist of type I collagen, whereas type II collagen comprises a large portion of the fibrillar collagen on the inner side (Cheung, 1987). Proteoglycans make for 2-3% of the dry weight and are mainly concentrated in the inner cartilaginous region of the meniscus (McDevitt and Webber, 1990; Buma et al., 2004). Also, the cellular component of the meniscus further reflects its fibrocartilaginous nature, the main cell type being meniscus fibrochondrocytes (McDevitt and Webber, 1990). Regarding cell types, at least two cell populations are present within the human meniscus (Ghadially et al., 1983). The fibrochondrocytes being the main cell type are reported within the inner and middle part of the meniscus having a rounded or oval shaped cell structure surrounded by an abundant ECM deposition (McDevitt and Webber, 1990; Ghadially et al., 1983). The outer one-third meniscus is reported to be populated mainly by spindle shaped fibroblast like cells with a dense connective tissue (Ghadially et al., 1983).\nOver the years, newer improvised methods such as meniscus allograft or autograft transplantation have been constantly searched for substituting the resected meniscus in case of either total or partial meniscectomy. However, none to date have generally been able to recapitulate and recreate the native meniscal multiporous and aligned structure as a single meniscus wedge shaped unit to completely and/or partially eliminate cartilage regeneration. As such, there is still a strong need to develop a scaffold that can mimic heterogeneous architecture and functions of native meniscal tissue."}
{"text": "1. Field of the Invention\nThe present invention relates to an EL (electro luminescence) display device having a semiconductor element (an element using a semiconductor thin film) formed on a substrate, and electric equipment using the EL display device as a display (display portion). The EL (electroluminescent) devices referred to in this specification include triplet-based light emission devices and/or singlet-based light emission device, for example.\n2. Description of the Related Art\nIn recent years, a technique by which a TFT is formed on a substrate has greatly progressed, and the application of that technique to an active matrix display device has been increasingly developed. In particular, a TFT using a polysilicon film enables high-speed operation since it is higher in field effect mobility (also called mobility) than a conventional TFT using an amorphous silicon film.\nAttention has been paid to the above active matrix display device because various advantages such as a reduction in manufacturing costs, a downsizing of the display device, an improvement in yield, and a reduction in through-put, are obtained by forming various circuits and elements on the same substrate.\nThe active matrix EL display device provides a switching element formed of a TFT (hereinafter referred to as a switching element) on each of pixels and activates a driver element that conducts a current control by the switching TFT, to thereby make an EL layer (speaking rigidly, a light emitting layer) emit light. For example, Japanese Patent Application Laid-open No. Hei 10-189252 discloses the EL display device.\nAs the active matrix EL display device, there are proposed two EL element structures depending on light radiating directions. One of those structures is that a light emitted from the EL element penetrates an opposing substrate and is then radiated so as to enter eyes of an observer. In this case, the observer can recognize an image from the opposing substrate side. The other structure is that a light emitted from the EL element penetrates an element substrate and is then radiated so as to enter eyes of the observer. In this case, the observer can recognize an image from the element substrate side.\nIn the former structure, the light from the outside penetrates the opposing substrate and is then irradiated onto the TFTs existing in gaps between the respective pixel electrodes, to thereby deteriorate the TFTs. However, because the light from the outside is not high in intensity, the deterioration of the TFTs is not large.\nOn the other hand, in the latter structure generally frequently employed, because the light emitted from the EL element penetrates the element substrate and is then radiated, the light emitted from the EL element is irradiated onto the TFTs, resulting in such a serious problem in that the TFTs are deteriorated.\nAlso, a storage capacitor is provided in the pixel and a high aperture ratio is demanded for the pixel from the viewpoint of the display performance. If the respective pixels have the high aperture ratio, the light application efficiency is improved, thereby being capable of achieving the power saving and the downsizing of the display device.\nIn recent years, the fine pixel size is developed, and a higher definition image is demanded. The fine pixel size increases an area of one pixel on which the TFT and the wirings are formed, to thereby reduce the pixel aperture ratio.\nUnder the above circumstances, in order to obtain the high aperture ratio of each pixel within the limit of a regular pixel size, it is essential to efficiently layout circuit elements necessary for the circuit structure of the pixel.\nAs described above, in order to realize the active matrix EL display device high in pixel aperture ratio with a small number of masks, an entirely novel pixel structure that has not existed up to now is demanded.\nThe present invention has been made to meet the above demands, and therefore an object of the present invention is to provide an EL display device having a pixel structure that realizes a high aperture ratio without increasing the number of masks and the number of processes.\nIn order to solve the problems with the conventional art, the present invention provides the following means.\nThe present invention is characterized by a pixel structure in which gaps between respective TFTs and gaps between respective pixels are shielded from a light without using a black mask. As one means for shielding the TFTs from the light, a gate electrode and source wirings are formed on a first insulating film, and most of a semiconductor layer that serves as an active layer is covered with gate wirings formed on a second insulating film different from the first insulating film. Also, as one means for shielding the gaps between the respective pixels from the light, pixel electrodes are so disposed as to be superimposed on the source wirings.\nThe above-mentioned TFTs are directed to switching TFT disposed on the respective pixels or current control TFTs.\nAccording to the structure of the present invention disclosed in this specification, there is provided an electronic device comprising a plurality of source wirings, a plurality of gate wirings, a plurality of current supply lines and a plurality of pixels, characterized in that:\neach of the plurality of pixels includes a switching TFT, a current control TFT, and a light-emitting element; and\nthe switching TFT includes a semiconductor layer (first semiconductor layer 200) having a source region and a drain region on an insulating surface, and a channel-forming region interposed between the source region and the drain region; a first insulating film (gate insulating film) formed on the semiconductor layer (first semiconductor layer 200); an electrode formed (first electrode 113) on the first insulating film so as to be superimposed on the channel-forming region; a source wiring (115) formed on the first insulating film; a second insulating film that covers the electrode (first electrode 113) and the source wirings; and a gate wiring (145) formed on the second insulating film and connected to the electrode (first electrode 113).\nIn the above structure, the electronic device is characterized in that the semiconductor layer (first semiconductor layer 200) has a region, which is superimposed on the gate wiring.\nFurther, the electronic device is characterized in that the region of the semiconductor layer which is superimposed on the gate wiring includes at least the channel-forming region, a region existing between the channel-forming region and the drain region, or a region existing between the channel-forming region and the source region, and is protected from light from the outside.\nIn case of the electronic device of a multi-gate structure in which a plurality of gate electrodes are on one semiconductor layer through an insulating film, it is characterized in that the semiconductor layer includes a plurality of channel-forming regions that the gate wiring is so disposed as to be superimposed on a region existing between one of the channel-forming regions and another channel-forming region.\nFurther, the electronic device is characterized in that the electrode and the source wirings are made of the same material on the first insulating film and that the pixel electrode, the connection electrode and the gate wiring are made of the same material on the second insulating film.\nAccording to another structure of the present invention, there is provided an electronic device comprising a plurality of source wirings, a plurality of first gate wirings, a plurality of current supply lines, a plurality of second gate wirings and a plurality of pixels, characterized in that:\neach of the plurality of pixels includes a switching TFT, a current control TFT, an erasing TFT and a light-emitting element; and\nthe switching TFT includes a semiconductor layer (first semiconductor layer 900) having a source region and a drain region formed on an insulating surface, and a channel-forming region interposed between the source region and the drain region; a first insulating film (gate insulating film) formed on the semiconductor layer (first semiconductor layer 900); an electrode (first electrode 805) formed on the first insulating film and superimposed on the channel-forming region; a source wiring (803) formed on the first insulating film; a second insulating film that covers the electrode (first electrode 805) and the source wiring (803); and a first gate wiring (801) formed on the second insulating film and connected to the electrode (first electrode 805).\nFurther, according to another structure of the present invention, there is provided an electronic device comprising a plurality of source wirings, a plurality of first gate wirings, a plurality of current supply lines, a plurality of second gate wirings and a plurality of pixels, characterized in that:\neach of the plurality of pixels includes a switching TFT, a current control TFT, an erasing TFT, and a light-emitting element; and\nthe erasing TFT includes a semiconductor layer having a source region and a drain region formed on an insulating surface, and a channel-forming region interposed between the source region and the drain region; a first insulating film (gate insulating film) formed on the semiconductor layer; a first electrode (third electrode 807) formed on the first insulating film and superimposed on the channel-forming region; a second electrode (second electrode 806) formed on the first insulating film; a second insulating film that covers the first electrode (third electrode 807) and the second electrode (second electrode 806); and a second gate wiring (802) formed on the second insulating film and connected to the first electrode (third electrode 807).\nIn the above structure, the electronic device is characterized in that the semiconductor layer has a region which is superimposed on the second gate wiring (802) and that the second gate wiring (802) is superimposed on at least the channel-forming region.\nFurther, the electronic device is characterized in that the region of the semiconductor layer which is superimposed on the second gate wiring (802) includes at least the channel-forming region, a region existing between the channel-forming region and the drain region, or a region existing between the channel-forming region and the source region, and is protected from light from the outside.\nIn the above structure, the first electrode (third electrode 807) that is superimposed on the channel-forming region comprises a gate electrode of the erasing TFT.\nIn the above structure, the second electrode (second electrode 806) comprises a gate electrode of the current control TFT, which is connected to the drain region of the switching TFT.\nFurther, it is characterized in that the first gate wiring and the second gate wiring are made of the same material in order to suppress the increase in the number of masks."}
{"text": "This invention relates, in general, to a dual language electronic reference machine. More particularly, this invention relates to a bilingual language teaching machine is which words in either of the two languages involved can be entered and information, in either language, is provided about the entered word. The information provided includes inflections of the translation of the input word and information about those inflections.\nA wide variety of electronic language reference products are available or taught in the literature. These products provide information concerning spelling, grammar and pronunciation. A number of products have been devised which disclose some dual language capability.\nU.S. Pat. No. 4,809,192 issued in February 1989 to Washizuka et. al. discloses a language translator with a speech synthesizer which in effect is a talking phrase book with ability to input certain words.\nU.S. Pat. No. 4,393,462 issued Jul. 12, 1983 to Tanimoto et. al. discloses a precision electronic translator that provides pronunciation of an input word and means for pronouncing the translated word equivalent to the input word.\nU.S. Pat. No. 4,489,396 issued Dec. 18, 1984 to Hashimoto et. al. is a translator with a voice synthesizer for pronunciation of the specific words in the input word language.\nThese and other translation devices provide translation of a word, and in some cases a phrase, together with pronunciation of the input word and/or pronunciation of the translated word or phrase. They sometimes also provide definitional information. They are useful in the limited context of a traveler. But they provide little or no introduction to the use of a language and are inadequate as a teaching or learning tool for learning about a second language.\nWhat is needed is a technique for providing in depth information about a foreign language to enable a user to employ the language in a more usual conversational fashion. This requires the presentation of not only a translation and the usual dictionary information but also the inflected forms of the word in the language being learned and a discussion of the inflected forms. These inflected forms include tenses and moods of verbs, masculine and feminine forms of the words and plurals. Furthermore, instructional material concerning the significance and meaning of the various inflected forms is important in order to understand the context of the usage.\nProviding all this information in a usable and useful manner for the person learning the language requires a mode of presentation that avoids burying the user in the information yet provides information necessary to enable the user to learn a language in depth. The mode of presentation is the key to the utility of an electronic language learning machine.\nAccordingly, a major purpose in this invention is to provide a dual language reference product that is optimally accommodated to the learning requirements of a user so as to facilitate the learning of a foreign language.\nIt is a related purpose of this invention to provide such a device as can be used by individuals having a wide range of competence in the language being learned.\nMore particularly, it is a purpose of this invention to provide a compact, hand held, self-contained dual language reference product that meets the above objectives."}
{"text": "The endothelial cell participates in numerous functions of vascular physiology. Many factors, such as cytokines, can alter the surface of the endothelial cell and thereby modulate the role of the endothelium in coagulation, inflammation, vaso-regulation, and adhesion. See, for example, R. P. Hebbel et al., J. Lab. Clin. Med., 129, 288 (1997); J. S. Pober, Am. J. Path., 133, 426 (1988); E. J. Favaloro, Inmmunol. Cell. Biol., 71, 571 (1993). The endothelial cell may also have a key role in the vascular pathology of sickle cell anemia, including the vaso-occlusions that cause acute painful crises. However, research in this area has been hindered by the inaccessibility of vascular endothelium in patients. For example, E. M. Levine et al. (U.S. Pat. No. 5,132,223) disclosed cloning and serial cultivation of adult human endothelial cells derived from brain-dead, but heart-beating cadaver organs. K. Gupta et al., Exp. Cell. Res., 230, 244 (1997) reported the culture of microvascular endothelial cells derived from newborn human foreskin. Thus, circulating endothelial cells might provide useful material for the study of vascular pathologies, for gene therapy, and for biomedical engineering applications. In previous investigations increased numbers of circulating endothelial cells have been found in sickle cell anemia and other conditions associated with vascular injury, such as that due to cytomegalovirus infection, rickettsial infection, myocardial infarction, intravascular instrumentation, and endotoxinemia. See, for example, F. George et al., Blood, 80, Suppl: 12a, abstract (1992); E. Percivalle et al., J. Clin. Invest., 92, 663 (1993), F. George et al., Blood, 82, 2109 (1993); C. A. Bouvier et al., Thomb. Diath. Haemorrh. Suppl., 40, 163 (1970); F. George et al., J. Immunol. Meth., 139, 65 (1991) and R. G. Gerrity et al., Exp. Mol. Pathol., 24, 59 (1976).\nHowever, in normal donors, there are only about 2-3 circulating endothelial cells per ml of blood; they have a quiescent phenotype, and about 50% of them are microvascular as evidenced by CD36 positivity. See, A. Solovey et al., New Engl. J. Med., 337, 1584 (1997), who reported using the methodology of Gupta et al., cited above, to coculture viable circulating endothelial cells identified in the blood of patients with sickle cell anemia with primary microvascular endothelial cells (MVEC). T. Asahara et al., Science, 275, 964 (1997) isolated putative endothelial cell (EC) progenitors from human peripheral blood after CD34+ enrichment by magnetic bead selection on the basis of cell surface antigen expression. The cells were cultured on fibronectin-coated wells in modified M-199 medium containing bovine brain extract and 20% fetal bovine serum. Q. Shi et al., Blood, 92, 362 (1998) characterized bone marrow-derived precursor endothelial cells by isolating CD34+ cells derived from peripheral blood using murine anti-CD34+ antibody binding followed by exposure to anti-mouse immunomagnetic beads. The cells were cultured in gelatin or fibronectin-coated plastic wells in M199 medium containing VEGF, FBS, bFGF and IGF.\nHowever, due to the low concentration of CEC in blood, a need exists for a culture method and medium that will permit the rapid expansion of CEC from blood, without the attendant difficulties of isolation discussed above."}
{"text": "1. Technical Field\nThe present invention relates to an improved data processing system and in particular to a method and apparatus for transmitting data. Still more particularly, the present invention relates to a method and apparatus for managing transmission of data from a source to a destination.\n2. Description of the Related Art\nTwo basic types of communications connections are employed between processors and between a processor and a peripheral. These types of connections are known as channels and networks. A channel provides a direct or switched point-to-point connection between communicating devices. This type of connection is typically employed between a processor and a peripheral device. The primary task of the channel is to transport data at the highest possible speed with the least delay. In contrast, a network is an aggregation of distributed nodes, such as workstations, file servers, and peripherals. Typically, in a network a node contends for the transmission medium and each node must be kept free of error conditions on the network. A traditional channel is hardware intensive and typically has lower overhead than a network. Conversely, networks tend to have relatively high overhead because they are software intensive. Networks, however, are expected to handle a more extensive range of tasks as compared to channels. In a closed system, every device addressed is known to the operating system either by assignment or pre-definition. This configuration knowledge is important to the performance levels of channels. Fibre Channel is a channel-network hybrid containing network features to provide the needed connectivity, distance, and protocol multiplexing along with enough traditional channel features to retain simplicity, repeatable performance, and guaranteed delivery. Fibre Channel has an architecture that represents a true channel/network integration. Fibre Channel allows for an active intelligent interconnections scheme, called a fabric, to connect devices. A Fibre Channel port manages simple point-to-point connection between itself and the fabric. A xe2x80x9cportxe2x80x9d is a hardware entity on a xe2x80x9cnodexe2x80x9d with a node being a device connected to a network that is capable of communicating with other network devices. Transmission is isolated from control protocol. As a result, different topologies may be implemented. Fibre Channel supports both large and small data transfers.\nThe demand for flexible, high performance, fault-tolerant storage subsystems caused host adapter, disk storage, and high-capacity drive manufacturers to adopt Fibre Channel (FC) as a standard. This serial standard cuts cabling costs, increases data rates, and overcomes distance limitations commonly associated with a Small Computer System Interface (SCSI). Fibre Channel can carry SCSI protocols, and as a result offers an ideal upgrade for work stations, servers, and other systems requiring high availability and/or high bandwidth. Fibre Channel has become increasingly important as companies are seeking to provide faster and easier access to data for various clients. The Fibre Channel Standard (FCS) as adopted by the American National Standards Institute (ANSI), provides a low cost, high speed interconnect standard for workstations, mass storage devices, printers, and displays.\nCurrent Fibre Channel data transfer rates exceed 100 megabytes (Mbytes) per second in each direction. Fibre Channel data transfer rates also may be scaled to lower speed, such as 50 Mbytes per second and 25 Mbytes per second. This technology provides an interface that supports both channel and network connections for both switched and shared mediums. Fibre Channel simplifies device interconnections and reduces hardware cost because each device requires only a single Fibre Channel port for both channel and network interfaces. Network, port to port, and peripheral interfaces can be accessed though the same hardware connection with the transfer of data of any format.\nIn sending data from a source node to a destination node, the source transmits data from a bus, such as a Peripheral Component Interconnect (PCI) bus, to a buffer for transfer onto a Fibre Channel system, which is connected to the destination node. Data is sent serially on Fibre Channel systems. As a result, data currently in a buffer must be sent before additional data may be loaded. Currently, if data cannot be sent because the destination is not accepting additional data, then this data must be removed to send data to another destination. This loading and dumping of data increases the overhead in transferring data between various nodes on a Fibre Channel system. Thus, it would be advantageous to have an improved method and apparatus for transferring data between nodes in which the overhead of acquiring the dumping and reloading of new data is eliminated.\nThe present invention provides a method and apparatus for transmitting data in a node having a buffer. A first set of data is received in a buffer for transmission to a target node. The first set of data is sent to the target node. Responsive to an indication that the target node is unable to receive data, a second set of data is loaded into the buffer for transmission to another target node, while the first set of data is retained in the buffer."}
{"text": "The Field of the Invention\nThe present invention relates to a digital broadcasting system for transmitting and receiving a digital broadcast signal, and more particularly, to a transmitting system for processing and transmitting the digital broadcast signal, and a method of processing data in the transmitting system and the receiving system.\nDescription of the Related Art\nThe Vestigial Sideband (VSB) transmission mode, which is adopted as the standard for digital broadcasting in North America and the Republic of Korea, is a system using a single carrier method. Therefore, the receiving performance of the digital broadcast receiving system may be deteriorated in a poor channel environment. Particularly, since resistance to changes in channels and noise is more highly required when using portable and/or mobile broadcast receivers, the receiving performance may be even more deteriorated when transmitting mobile service data by the VSB transmission mode."}
{"text": "The present invention relates to a table for holding, by suction under vacuum, an original document to be imaged by an imaging apparatus or the like. More particularly, the invention relates to an original table suitable for use with an imaging apparatus that performs imaging of an original to provide an image on a printing plate or for use with a platemaking apparatus used for this purpose.\nA typical configuration of an original table used with a conventional imager for plate making or platemaking apparatus is shown in FIG. 1. Specifically, an original mount 1 equipped at one side with a pivotable glass holder 2 holds the original 3 in position by means of the glass 2 during imaging.\nA problem with the original table using a holder glass is that the glass is heavy and makes the handling of the table cumbersome. Furthermore, at low reduction ratios, the original table must be placed only a short distance from the optical system, and this prevents the glass holder from opening sufficiently to achieve perfect registry of the original."}
{"text": "As is well known to those skilled in the art, various catalysts are used in processing. Many of these catalysts are characterized by the presence of catalytically active components on a support. Attempts are constantly being made to improve the properties of the support and to thus permit attainment of a catalyst composition, containing support preferably plus other ingredients, which is characterized by desirable properties including, for example, conversion, yield, selectivity, etc.\nIt is an object of this invention to provide a novel alumina, and a process for making this product. Other objects will be apparent to those skilled in the art."}
{"text": "1. Technical Field\nThe relates to a phase-change random access memory (PCRAM) device, and more particularly, to a PCRAM device and a method of manufacturing the same.\n2. Related Art\nWith demands on lower power consumption, next-generation memory devices having nonvolatile and non-refresh properties have been studied. A PCRAM device of the next-generation memory devices includes a switching element connected at intersections of word lines and bit lines, which are arranged to cross each other, a lower electrode electrically connected to the switching element, a phase-change layer formed on the lower electrode, and an upper electrode formed on the phase-change layer.\nIn a conventional PCRAM device, when a write current flows through the switching element and the lower electrode, Joule heat is generated at an interface between the phase-change layer and the lower electrode. The phase-change layer is phase-changed into an amorphous state or a crystalline state by the generated joule heat. Therefore, the conventional PCRAM device stores data using a difference between resistances in the amorphous state and the crystalline state of the phase-change layer.\nHowever, in the conventional PCRAM device, the Joule heat generated when the write current flows affects a phase-change layer of adjacent cell.\nThe effect on adjacent cells is generally referred to as thermal disturbance. In recent years, the thermal disturbance has an increased effect on adjacent cells when a semiconductor memory device is highly integrated.\nFIGS. 1A and 1B are views illustrating thermal disturbance of a conventional PCRAM device.\nAs shown in FIGS. 1A and 1B, the conventional PCRAM device includes a lower electrode 10 formed on a switching element (not shown), a phase-change layer 20 formed on the lower electrode 10, and an upper electrode 30 formed on the phase-change layer 20. The reference numeral 40 denotes an insulating layer.\nAs shown in FIG. 1A, if a cell A is written when cells B are written with data “1”, which is a high resistance state, Joule heat is generated at an interface between the lower electrode 10 and the phase-change layer 20 of the cell A (see FIG. 1B), and thus, phase-change material patterns of amorphous states in the cells B are crystallized. Therefore, resistances of the cells B are reduced.\nThe thermal disturbance generated in the conventional PCRAM device may cause a malfunction, and thus reliability of the conventional PCRAM device is degraded."}
{"text": "1. Field of Invention\nThis invention relates generally to optical systems for the storage and retrieval of information and, more particularly, to the read/write head of the magneto-optical information storage system which directs radiation to the storage medium and then directs radiation resulting from the interaction with the medium to the radiation detectors.\n2. Description of the Related Art\nThe optical storage systems, at present can generally be placed into one of two categories, the categories determined by the optical property used to identify different logical states on the storage medium. The first optical storage system can be referred to a differential absorption (or reflection) of a radiation beam impinging on the storage medium surface. In the differential absorption optical systems, each logical states are associated with changes in the intensity of a beam of radiation interacting with the storage medium. In the second category of optical storage systems, changes in the rotation of plane polarized beam of radiation are used to identify optical states. The present invention is directed to the second or the magneto-optical storage and retrieval systems and provides a technique for determining selected parameters in the optical path to enhance the identification of the two orientations of magnetic regions, orientations which encode the data stored on the disk.\nReferring to FIG. 1, the implementation of the read/write head in a magneto-optical information storage system, the system relying on differential rotation of the planar polarization of a optical radiation caused by the interaction of the optical radiation with the storage surface, is shown. This type of storage system relies on the Kerr effect wherein the rotation of a plane of polarization is different when a magnetic material has a magnetic orientation parallel to or a magnetic orientation anti-parallel to the direction of the radiation interacting with the magnetic material, i.e., the differential change in polarization of a reflected beam depends upon the orientation of the magnetization of the local domain with which the radiation interacts. As with the implementation for detecting a change in reflected light amplitude, the radiation from a light source 10 is collimated by lens 11 and one plane of polarization is selected by passing the collimated beam through the partial beam splitter 12'. Because linearly polarized radiation can be considered to be comprised of two circularly polarized radiation components, the interaction with the magnetic layer forming a portion of storage medium 15 effects the two circularly polarized components differently. As a result, after interaction with the storage material, the reflected radiation is not linearly polarized parallel to the applied radiation, but an elliptical polarization of the reflected radiation results in a rotation of the reflected linear polarization due to the circular dichroism and the circular birefringence of the storage media. The reflected radiation is recollimated by objective lens 14. The recollimated beam is applied to beam splitter 12 and the components of the radiation beam orthogonal to the plane of polarization of the radiation impinging on storage medium, i.e., the components induced by the interaction, are reflected by the beam splitter 12. Some of the light with polarization parallel to the impinging radiation can also be reflected from the magneto-optic region. The radiation reflected by the beam splitter 12 is transmitted through a quarter wave plate 16A and a half wave plate 16B to correct for ellipticity introduced into the radiation beam. The polarization beam splitter 17 divides the radiation reflected from beam splitter 12 into radiation components which have been rotated by the interaction with the storage material. Each detector 18 and 19 receives a component resulting from one orientation of the magnetic regions of the storage medium interacting with the impinging radiation beam. The differential amplifier 20 is used to enhance the detectability of the small signals, the rotation due to the Kerr effect typically being less than 2.degree. relative to reflected radiation which had not been subjected to differential interaction of the circularly polarized components with the optical storage material and to cancel the large DC component of the two radiation components.\nIn the optical storage systems using a magneto-optical storage medium, a need has been felt for a technique of determining how to optimize the parameters of the system in order to achieve the most detectable signal. In the article by W. A. Challener and T. A. Rinehart, \"Jones Matrix Analysis of Magnetooptical Media and Read-Back Systems\", Appl. Opt. 26, 3974 (1987), part of the problem of a differential detection system was addressed. In that article, the substrate birefringence and the wave plate tolerances were studied. However, the DC offset in the differential signal was not considered and a range of \"ideal\" wave plates was found, each with a sensitivity to the optical path birefringence. Therefore, the need has remained for generally applicable technique for identifying the parameters which would permit optimization of the detection of the state of the region of the storage system to which radiation was being applied."}
{"text": "1. Field of the Invention\nThe present invention relates to a design system, a design method, and a storage medium storing a design program for a structural analysis after amending the form of a model.\n2. Description of the Related Art\nConventionally, when a desirable model form is obtained to minimize the weight while maintaining the strength in analyzing a structure, a plurality of candidates for a desirable form (basis vector) are defined, and then the form (a) is divided into mesh units as shown in FIG. 1 into the form (b). Then, in a manual operation, as indicated by the form (c), the coordinate indicated by a small circle (∘) of a node on each of the edges (sides) of the divided mesh units is manually input (by specifying the x, y, and Z coordinates) and amended. A structural analysis is made on a model formed based on the amended edge according to the well-known boundary conditions and the attribute (strength, weight, etc.) of a material to obtain the optimum solution satisfying a predetermined strength, for example, a model form lightest in weight.\nWhen the above described optimization is performed, it is necessary to divide a model into mesh units, and manually associate a node on each edge (side) with a coordinate value as shown by (c) in FIG. 1. If the number of nodes is large, then a great number of manual operations are required, thereby preventing a quick structural analysis.\nTo solve the above described problems, the present invention aims at quickly performing a structural analysis in a simple operation by dividing a model into mesh units, specifying a child edge of an optional curve in a parent edge, projecting the parent edge into the child edge, automatically computing the coordinates, and obtaining the optimum solution in a structural analysis based on the computation.\nAccording to the first aspect of the present invention, the design system which analyzes a structure by amending the form of a model includes: a mesh division unit for dividing the model into mesh units; a child edge generation and display unit for generating and displaying a child edge, which is newer than a parent edge divided into mesh units; an amendment unit for amending the displayed child edge into an optional form, an association unit for obtaining the correspondence between the above described parent unit and the amended child unit; and a structural analysis unit for analyzing the structure of the amended model form depending on the obtained correspondence."}
{"text": "The invention relates to a method for decontaminating nuclear reactors employed to generate electric power and more particularly to a method for performing full system decontaminations on boiling water reactors.\nDuring on-line power generating operations of commercial boiling water nuclear reactors, thin layers of metal oxides tend to build up on the internal surfaces of vessels and other components and piping in contact with circulating primary coolant (essentially high temperature water). Activated metal ions in the central core regions in reactor pressure vessels are entrained in the primary coolant and then are absorbed in the metal oxides, which results in relatively high radiation levels on these surfaces. It is desirable to reduce the radiation levels to xe2x80x9cAs Low As Reasonably Achievablexe2x80x9d levels in order to reduce the exposure of personnel working near the reactors during periodic plant outages and/or plant decommissioning operations. Thus, the industry may employ one or a combination of various known chemical decontamination treatments, e.g., acid permanganate, alkaline permanganate, Citrox, CAN-DEREM, LOMI and/or other processes, in order to dissolve or break up the oxide films. Conventionally, these decontamination processes involve the addition of permanganate, oxalate, citrate, EDTA and/or other ions to the primary coolant to form decontamination solutions and then the circulation of the solutions through the components to be decontaminated. In addition to removing the oxide layers, it may be desirable to remove several microns of base metal in order to better protect personnel during decommissioning processes. Dilute chemical decontamination solutions generally contain less than about 3-5% by weight of such decontamination agents. Chemical decontaminations may be performed upon full primary coolant systems or upon selected subsystems. Full system decontaminations are the preferred approach when the goal is to reduce dose rates on multiple subsystems throughout the plants. In addition, full system decontamination processes are generally performed with nuclear fuel assemblies out of the central core regions of the reactor pressure vessels, but the fuel assemblies may be retained in the central core regions in some cases.\nThe activated metal ions that are removed from the internal surfaces of the primary coolant systems in the course of the decontamination operations are collected on cation exchange resins. The activated resins must then be removed to remote disposal sites.\nThe majority of the activated oxide deposits in boiling water reactor primary coolant systems are located in the central core regions of reactor pressure vessels. These deposits do not substantially contribute to personnel exposure. Thus, it would be very desirable to decontaminate only those systems that substantially contribute to personnel exposure and bypass the central core regions. This would substantially reduce the total exposure of personnel while reducing resin and disposal costs.\nIt is an object of the present invention to decontaminate a portion of a reactor pressure vessel in a boiling water reactor and its appurtentant recirculation system while bypassing its central core region. It is a further object to substantially decontaminate a boiling water reactor with lower overall personnel exposures to radiation and lower resin costs.\nWith these objects in view, the present invention resides in a method of decontaminating a boiling water reactor having a plurality of reactor recirculation loops hydraulically connected in parallel with a reactor pressure vessel. Such a reactor pressure vessel has: a central core region; an annulus region surrounding the central core region and in hydraulic communication with the recirculation loops; and a lower internals region in hydraulic communication with the central core region. In the practice of the present invention, a decontamination solution is circulated through at least one of the reactor recirculation loops and the annulus region of the pressure vessel without circulating through the central core region. In a preferred practice of the present invention, the decontamination solution also circulates between the annulus region and the lower internals region without circulating through the central core region. Thus, a boiling water reactor can be substantially decontaminated while reducing overall personnel exposure and generating less resin wastes."}
{"text": "Blood clot formation, or “thrombosis,” is a basis of a number of serious diseases, such as ischemic stroke, myocardial infarction (heart attack), and deep vein thrombosis (DVT). Blood clots, or “thrombi,” form inside blood vessels and obstruct the flow of blood through the circulatory system, thereby depriving tissue and organs of oxygen. In the case of a stroke, for instance, when blood flow to the brain is obstructed for longer than a few seconds, brain cells can die and permanent neurological damage can result.\nThrombi can be treated (reduced or eliminated) by inducing thrombolysis. Thrombolysis is the dissolving, or “lysis,” of a thrombus. Thrombolysis can sometimes be induced pharmacologically, such as by administering a tissue plasminogen activator drug (tPA), the most common thrombolytic agent. Thrombolytic agents (commonly called “clot-busting drugs”) can be administered via an intravenous line or using a catheter to deliver them proximally to the thrombus. However, thrombolysis by administration of clot-busting drugs has its limitations. For example, to be successful, the clot-busting drugs should be administered within three (3) hours of an acute ischemic stroke, and preferably within two (2) hours. Further, patients who use blood-thinning medications, and certain other medications, are usually not candidates for pharmacological thrombolysis. And of those patients receiving the treatment, it is unsuccessful in dissolving thrombi in approximately 25% of patients.\nIn view of the limitations of pharmacologically induced thrombolysis, various medical devices for surgically removing thrombi have been developed. The procedure for surgically removing thrombi is generally known as a “thrombectomy.” In thrombectomy treatments, a catheter system is typically used to deliver a device to the thrombus. The device can be, for example, an aspiration catheter. Aspiration catheters can perform a thrombectomy by suctioning the thrombus out of the blood vessel. Other thrombectomy procedures use a mechanical device to physically entangle with a thrombus, and to remove the thrombus as the device is removed from the blood vessel. Various types of mechanical devices, such as wires, corkscrew-like coils, bristles, and baskets have been employed to entangle with thrombi.\nSome traditional thrombectomy devices can cause damage to blood vessel walls. In addition, some traditional thrombectomy devices can be prone to generating thrombotic fragments that become emboli when they travel within the bloodstream. Emboli can become lodged in arteries, veins, arterioles, and capillaries, and can block the blood supply to vital organs such as the brain or heart. Emboli in the bloodstream can be life-threatening. In the case of DVT treatment, dislodged thromboemboli can travel to the lungs, resulting in a pulmonary embolism, which can be fatal."}
{"text": "The present invention relates to electrochemical conversion cells, commonly referred to as fuel cells, which produce electrical energy by processing first and second reactants. For example, electrical energy can be generated in a fuel cell through the reduction of an oxygen-containing gas and the oxidation of a hydrogenous gas. By way of illustration and not limitation, a typical cell comprises a membrane electrode assembly (MEA) positioned between a pair of flowfields accommodating respective ones of the reactants. More specifically, a cathode flowfield plate and an anode flowfield plate can be positioned on opposite sides of the MEA. The voltage provided by a single cell unit is typically too small for useful application so it is common to arrange a plurality of cells in a conductively coupled “stack” to increase the electrical output of the electrochemical conversion assembly.\nThe membrane electrode assembly typically comprises a proton exchange membrane separating an anode layer and a cathode layer of the MEA. The MEA is typically characterized by enhanced proton conductivity under wet conditions. For the purpose of describing the context of the present invention, it is noted that the general configuration and operation of fuel cells and fuel cell stacks is beyond the scope of the present invention. Rather, the present invention is directed to methods for managing MEA hydration cycling fatigue life in fuel cells. Regarding the general configuration and operation of fuel cells and fuel cell stacks, applicants refer to the vast collection of teachings covering the manner in which fuel cell “stacks” and the various components of the stack are configured. For example, a plurality of U.S. patents and published applications relate directly to fuel cell configurations and corresponding methods of operation. More specifically, FIGS. 1 and 2 of U.S. Patent Application Pub. No. 2005/0058864 and the accompanying text present a detailed illustration of the components of one type of fuel cell stack and this particular subject matter is expressly incorporated herein by reference."}
{"text": "This invention relates to the simultaneous measurement of the concentration of a selected ion species in a solution and the pH of the solution. The invention particularly, though not exclusively, relates to photographic solutions, and particularly, though not exclusively, where the selected ion species is silver. In general, however, the invention relates to the simultaneous potentiometric measurement of the concentration of any ion species in a solution and measurement of the pH of the solution using an ISFET (Ion Selective Field Effect Transistor).\nFor the present purpose the tern xe2x80x9csolutionxe2x80x9d is to be understood as also including an emulsion, for example a mixture of a silver compound suspended in gelatin, or a dispersion. The invention will be particularly described, by way of example only, with reference to photographic solutions.\nIt is known simultaneously to measure silver ion concentration in, and the pH of, an aqueous solution. In one arrangement, a single reference electrode is connected into a first potentiometer circuit with a conventional glass pH electrode, and is connected into a second potentiometer circuit with a conventional silver electrode, all three electrodes being immersed in the solution. In another arrangement, an ISFET is used instead of the glass pH electrode. This necessitates the use of a separate reference electrode for each measuring circuit in order to provide electrical isolation between the circuits since the ISFET is a current carrying device whose presence would otherwise interfere with the voltage measurement of the silver electrode.\nA glass pH electrode has the disadvantage that it can be damaged under conditions of high temperature and high pH, so that its readings become unreliable or inconsistent. An ISFET overcomes this disadvantage. However, the conventional arrangement including an ISFET described above is complicated by the requirement of the additional reference electrode, especially when applied in a large scale production vessel, as used in the preparation of photographic emulsions for example, where the electrodes are configured in a unitary probe. This can lead to difficulties for maintenance and for calibration. Furthermore, existing probe structures would require extensive modification to accommodate the additional reference electrode, which would be expensive.\nIt will be appreciated that if, on the other hand, measurement of ion concentration and pH were not required simultaneously, then the measurements would not interfere with each other and a single reference electrode could be used successively in combination with an ion concentration electrode and an ISFET.\nIn accordance with one aspect of the present invention, there is provided apparatus for simultaneously measuring the concentration of a selected ion species in a solution and the pH of the solution, comprising: a first electrical circuit that is arranged to receive signals from both a reference electrode and an ion selective electrode immersed in the solution and to derive therefrom an output signal representative of the concentration of the selected ion in the solution; a second electrical circuit that is arranged to receive signals from both said reference electrode and an ISFET immersed in the solution and to derive therefrom an output signal representative of the pH of the solution; wherein any d.c. input signal to said first electrical circuit from the reference electrode is substantially electrically isolated from the input of the second circuit; wherein a signal representative of the voltage, usually earth potential, of the solution is supplied (a) directly to the first circuit so as to establish a reference, usually earth, potential for the first circuit, and (b) to the second circuit through a.c. coupling means so as to establish a corresponding virtual reference, usually earth, potential for the second circuit; and wherein the first and second electrical circuits are arranged to be provided with electrical power from supplies that are electrically isolated from each other.\nThe apparatus may comprise means for displaying a representation of said ion concentration and pH output signals, wherein said second electrical circuit includes an isolation amplifier, and wherein said display means is arranged to receive said pH output signal of the second circuit through the isolation amplifier. Preferably, the apparatus includes a further isolation amplifier through which the ion concentration output signal of the first circuit is supplied to the display means. Advantageously, the apparatus comprises a low pass filter, wherein said pH output signal from the second electrical circuit is arranged to be passed to the display means through the low pass filter.\nPreferably, the apparatus comprises a high value resistor, for example of about 1 Mxcexa9 or greater, that is arranged to effect said electrical isolation of d.c. input signals to said first and second electrical circuits. Also said a.c. coupling means may comprise a high value capacitor, for example of about 1 xcexcF or greater.\nIn accordance with another aspect of the present invention, there is provided a method of simultaneously measuring the concentration of a selected ion species in a solution and the pH of the solution, comprising the steps of: measuring in a first electrical circuit the potential difference between an ion selective electrode and a reference electrode both immersed in the solution, and deriving therefrom the concentration of the ions in the solution; measuring in a second electrical circuit the current flowing between an ISFET and the reference electrode both immersed in the solution, and deriving therefrom the pH of the solution; connecting the reference electrode to the first and second electrical circuits such that any d.c. signal from the reference electrode is electrically isolated from the second circuit; making an electrical connection between the solution and the first circuit so as to provide the solution potential as a reference, preferably earth, potential therefore, and making an electrical connection between the solution and the second circuit through a.c. coupling means so as to provide a corresponding virtual reference, preferably earth, potential therefore; and supplying the first and second circuits with electrical power from sources that are electrically isolated from each other.\nThe method of the invention is advantageously carried out using the apparatus of the invention.\nDetails of electrodes suitable for use in the present invention as ion selective and reference electrodes, and of ISFETs, can be found in the book xe2x80x9cpH Measurementxe2x80x9d by Helmuth Galster (VCH,1991).\nThe electrical isolation of the two circuits provided in the present invention allows an ISFET to be used in the pH measuring circuit, whilst needing only a single, common, reference electrode. The disadvantages of the known arrangements for simultaneous ion concentration and pH measurement are thus overcome in a particularly convenient manner.\nThe isolation is provided at several stages. Initially this is done by arranging that the signal from the reference electrode is used in the ion concentration circuit as a potentiometric measurement, and is supplied to the pH measuring circuit only as an a.c. input, i.e. after having any d.c. component isolated therefrom. An actual reference potential, the potential, usually earth, of the solution, is applied to the first circuit, and a virtual reference potential derived therefrom is applied to the second circuit. The two circuits have separate isolated power supplies. Furthermore, when the resulting ion concentration and pH signals are supplied to a display means, such as a multi-channel voltmeter, this is done through respective isolation amplifiers, which are preferably supplied from a third, isolated power supply.\nThe ability to use a single reference electrode means that a single, unitary measurement probe can be constructed, in which the ISFET can be installed relatively easily along with the ion selective and reference electrodes. Furthermore, the measuring apparatus can be calibrated more easily than is the case with the known arrangement using two reference electrodes.\nAlthough described with reference to a single ion selective electrode and a single ISFET, it is envisaged that the present invention may comprise two or more ion concentration electrodes and/or two or more ISFETs, each type of electrode being connected into the respective first or second electrical circuit.\nApparatus for, and a method of, simultaneously measuring the concentration of a selected ion species in a solution and the pH of the solution, will now be described, by way of example, with reference to the accompanying schematic circuit diagram."}
{"text": "The present application relates generally to methods and systems for making payments for recyclable items such as beverage containers deposited at reverse vending machines.\nContainer deposit laws also known as “bottle bills” require consumers to pay a small refundable deposit on containers (such as beer, soft drink, and other beverage containers) they purchase to improve the rate of recycling or reuse.\nBottle bills in most states typically work as follows. When a retailer buys beverages from a distributor, the retailer pays a deposit to the distributor for each can or bottle purchased. The consumer pays the deposit to the retailer when buying the beverage. When the consumer returns the empty beverage container to the retail store, a redemption center, or a reverse vending machine, the deposit is refunded. The retail store or redemption center recoups the deposit from the distributor, typically along with an additional handling fee to help cover the cost of handling the containers.\nAs used herein, the term “distributor” refers to a product maker (e.g., the company producing the product such as the Coca-Cola Company) or to an entity distributing products from a product maker to retailers.\nReverse vending machines are automated machines that accept used beverage containers and return container deposits to users. Manufacturers of reverse vending machines include Tomra, Wincor Nixdorf, Envipco, Repant, reVend, and Can and Bottle Systems.\nIn use, the user places the empty beverage container into a receiving aperture of the reverse vending machine. The reverse vending machine automatically identifies the beverage container, e.g., by scanning the beverage container's UPC code. Eligible containers are accepted, and the user is paid. The user is typically given a receipt that can be redeemed at a retailer checkout station. Accepted containers are typically crushed or broken in the reverse vending machines to reduce their size."}
{"text": "Carousels, merry-go-rounds, playground rides, or playground roundabouts have long been a popular equipment on playgrounds for children of all ages, and can be found in numerous different configurations, sizes, functionality, and designs.\nTypically such playground roundabouts comprise some sort of platform or seats to accommodate the one or more users and which then is turned either mechanically, by others or by the persons on the roundabout themselves.\nDifferent types of roundabouts exist yielding different types of movement of the user not only making the user circle in a horizontal plane but also move up and down, or spin around his own axis etc.\nU.S. Pat. No. 3,073,595 describes a playground ride where the seating platform is arranged on two freely rotating and separated hubs whose axis of ration are angularly offset resulting in a special combination of circular motion and vertical motion. As the seating platform can be rotated fast or slowly and in changing direction relative to the rotation of the base column is obtained a more exhilarated ride where the resulting movement is perceived as unpredictable by the user.\nAnother type of exercising device is described in U.S. Pat. No. 4,290,601 comprising a wobble plate mounted on a single shaft to a base and placed on one side on top of a roller assembly thereby creating a raised position of one portion of the wobble plate. The roller assembly is in one embodiment made from pair of hemispherical balls and held in its radial position by a radial arm connected to the shaft. A person standing on the wobble plate can then set the board in motion. The design with the roller assembly moving between the base and the underside of the wobble board however yields a construction with a relatively high number of movable parts and with a number of places where objects such as clothes, fingers etc. can get caught or entangled. The latter especially makes the exercise device unsuitable as a playground equipment.\nEP 1747803 discloses an exercising apparatus of an oscillating platform placed at an angle relative to a support platform and free to rotate around two non-parallel axes. Different configurations of the plate and the axes relative to each other yield different types of motion of the user. Also, the platform rotating assembly may be driven by a gear assembly configured to rotate the platform at a predetermined rate such that the platform rotates relative to the platform by a predetermined amount during use. Thereby the amount of oscillating can be controlled however without influencing the type of motion of the user."}
{"text": "The present invention relates to an erasing method of a floating gate type nonvolatile semiconductor storage device.\nIn recent years, there has been a demand for reducing the power consumption in accordance with increase of integration level in a flash memory of a nonvolatile semiconductor storage device. In response to the above demand, a reduction in consumption of power is enabled by using the Fowler-Nordheim tunneling phenomenon for write (program) and erase operations. The flash memory that executes the write and erase operations utilizing the Fowler-Nordheim (referred to as FN hereinafter) tunneling phenomenon is called the FN--FN type flash memory.\nOn the other hand, flash memories are classified by the memory cell array structure, and four principal types will be enumerated hereinbelow.\n[1] The Institute of Electronics, Information and Communication Engineers Technical Report, ICD93-128, p37, 1993\nAn AND type flash memory reported as \"\"AND\" cell structure for a 3V-only 64Mbit Flash Memory\"\n[2] The Institute of Electronics, Information and Communication Engineers Technical Report, ICD93-26, p15, 1993\nA DINOR type flash memory reported as \"A Novel Cell Structure Suitable for a 3 Volt Operation, Sector Erase Flash Memory\"\n[3] IEDM Technical Digest, p263-266, 1995\nA DuSNOR type flash memory reported as \"A Novel Dual String NOR (DuSNOR) Memory Cell Technology Scalable to the 256 Mbit and 1 Gbit Flash Memories\"\n[4] IEDM Technical Digest, p267-270, 1995\nAn ACT type flash memory published in \"A New Cell Structure for Sub-quarter Micron High Density Flash Memory\" and \"A sensing Scheme for a ACT flash memory\" of The Institute of Electronics, Information and Communication Engineers Technical Report, ICD97-21, p37, 1997\nThe above types are published by several companies.\nAccording to the flash memories of the above types [1] through [4], it is acceptable to execute electrical writing (program) and erasing on a memory cell. However, a voltage is applied to the drain, source or control gate of the selected memory cell in the write operation and the erase operation, while a voltage is also applied to the drain, source or control gate of the unselected memory cell. The threshold voltage of the unselected memory cell is changed by the influence of this voltage application, possibly causing erroneous reading\nIn recent years, there is an increasing trend toward using a method for applying a negative voltage to the semiconductor substrate (well) in order to reduce the absolute value of a voltage to be used inside the flash memory in the erase operation. This negative voltage applied to the semiconductor substrate brings the unselected memory cell whose drain, source or control gate receives the voltage into a lightly erased state, exerting bad influence (referred to as a substrate disturbance hereinafter) on the threshold voltage of the unselected memory cell. The substrate disturbance tends to become more severe as the flash memory comes to have a larger capacity.\nThe aforementioned substrate disturbance will be described by taking the ACT (Asymmetrical Contactless Transistor) type flash memory as an example.\nFIG. 6 shows a sectional view of one memory cell of the above ACT type flash memory, and the principle of operation of the ACT type flash memory will be described with reference to FIG. 6.\nIn the above ACT type flash memory of FIG. 6, a tunnel oxide film 14, a floating gate 15, an interlayer insulating film 16 and a control gate 17 are lamellarly formed on a substrate (p-type well) 11 so as to form a bridge between a drain 13 and a source 12 formed on the substrate 11. Then, the drain 13 and the source 12 have different donor concentrations.\nIn the case of a program operation in the ACT type flash memory having the aforementioned construction, that is, in the case where electrons are extracted from the floating gate 15 to provide a written state (data \"0\"), a negative voltage Vnw (-8 V) is applied to the control gate 17 and a positive voltage Vpp (+5 V) is applied to the drain 13, thereby extracting electrons from the floating gate 15 by the Fowler-Nordheim (referred to as FN hereinafter) tunneling phenomenon with the source 12 brought into the floating state. By this a program operation, the threshold voltage of the memory cell is lowered to a voltage of about 1.5 V.\nIn the case of an erase operation, that is, in the case where electrons are injected into the floating gate 15 to provide an erased state (data \"1\"), a positive voltage Vpe (+10 V) is applied to the control gate 17, a negative voltage Vns (-8 V) is applied to the source 12, and the drain 13 is brought into the floating state. Electrons are injected into the floating gate 15 by the FN tunneling phenomenon. Therefore, the threshold voltage of the memory cell is increased to a voltage of about 4 V or more.\nIn the case of a read operation, a voltage of +3 V is applied to the control gate 17, a voltage of +1 V is applied to the drain 13, and a voltage of 0 V is applied to the source 12. The data is read by the sensing circuit (not shown) for sensing the current flowing through the memory cell.\nThe voltages applied to the memory cell in the program, aforementioned operations are shown in Table 1.\nTABLE 1 Control Substrate Gate Drain Source P-Type Well Program -8 V 5 V Open 0 V Operation Erase 10 V Open -8 V -8 V Operation Read 3 V 1 V 0 V 0 V Operation\nIn order to explain the substrate disturbance in the erase operation, the application voltage in the erase operation will be described with reference to the array structure of the ACT type flash memory shown in FIG. 7. As schematically shown in FIG. 7, the array structure of the ACT type flash memory has a virtual-ground-type array structure in which two memory cells jointly own an identical bit line.\nIn the above ACT type flash memory are shown main bit lines BL0 through BL4096, sub-bit lines SBL00 through SBL04096 and SBL10 through SBL14096 formed from a diffusion layer (the sub-bit lines being in a layer different from that of the main bit lines), word lines WL0 through WL63, selection gate signal lines SG0 and SG1 of selection transistors ST00 through ST04096 for selecting each block and a contact section CN (the portions each being indicated by the mark .box-solid. in FIG. 7) of the main bit lines BL0 through BL4096 and the sub-bit lines SBL00 through SBL04096 and SBL10 through SBL14096. Then, in regard to the memory cells M00, M01, ... , M10, M11, ..., the number of contacts is reduced by making the memory cells of adjoining lines jointly own the sub-bit lines SBL01 through SBL04095 and SBL11 through SBL14095 and using the diffusion layer for the sub-bit lines SBL00 through SBL04096 and SBL10 through SBL14096, by which the array area is sharply reduced, allowing high-density integration to be achieved.\nFIG. 8 schematically shows the sub-bit lines SBL00 through SBL04096 and SBL10 through SBL14096 (shown in FIG. 7) formed from the aforementioned diffusion layer in the form of a cross-section of the essential part of the ACT type flash memory.\nAs shown in FIG. 8, an interlayer insulating film 22, a floating gate 23 (FG) and a control gate 24 (WL) are lamellarly arranged on a semiconductor substrate 20 on which a sub-bit line 21 (diffusion layer) is formed. Then, the common sub-bit line 21 provided below the end portion of adjoining floating gates 23 (FG) has donor concentrations that differ between a drain 21a and a source 21b.\nIn the case of the aforementioned ACT type flash memory, the erasing operation is executed on a block basis. In the erase operation, for example, a positive voltage (+10 V) is applied to the word lines WL0 through WL31 connected to the control gates of the memory cells M00, M01, . . . of a selected block BLOCK0 shown in FIG. 7 in order to increase the threshold voltage of the memory cells. Further, a negative voltage (-8 V) is applied to a semiconductor substrate (well) and the main bit lines BL0 through BL4096. In this stage, the selection gate signal line SG0 has a voltage of 0 V to turn or the selection transistors ST00 through ST04096, and a negative voltage (-8 V) is applied to the sub-bit lines SBL01 through SBL04095. By this operation, a high electric field is generated between the floating gates and the channels of the memory cells M00, M01, . . . , by which electrons are injected into the floating gate by the FN tunneling phenomenon, increasing the threshold voltage of the n memory cells M00, M01, . . . to a voltage of 4 V.\nOn the other hand, in an unselected block BLCK1 in FIG. 7, a reference voltage Vss (0 V) is applied to the word lines WL32 through WL63. When a negative voltage (-8 V) is applied to the selection gate signal line SG1, then the selection transistors ST10 through ST14096 are turned off, as a consequence of which the sub-bit lines SBL10 through SBL14096 connected to the selection transistors ST10 through ST14096 are brought into the floating state. In this stage, the semiconductor substrate is common to all the memory cells. Therefore, the negative voltage (-8 V) is applied to the substrate, and an electric field is generated between the floating gate and the semiconductor substrate although the above electric field is less than that of the foregoing selected block. This electric field causes injection of electrons into the floating gate. The injection of electrons into the floating gate in the unselected block more frequently occurs in the memory cell in the low threshold voltage state, i.e., in the memory cell in the programmed state, i.e., in the memory cell of data \"0\".\nHere is now considered the substrate disturbance in a 64-M flash memory in which 512 blocks each having a block size of 16 KB exist. If each block has been subjected to one million times of rewriting, assuming that each erasing time is 2 ms, then a disturbance time obtained by summing up the times applied to the unselected block in the above case is expressed by: EQU 511.times.1,000,000.times.2 msec .apprxeq.10.sup.6 sec (1)\nFIG. 9 shows an example of the substrate disturbance in the erase operation. In FIG. 9, the horizontal axis represents the disturbance time, while the vertical axis represents the threshold voltage Vt (conditions:control gate voltage Vg of 0 V; drain Vd and source voltage Vs of floating; and substrate voltage Vsub of -8 V). As is apparent from FIG. 9, the threshold voltage of the memory cell becomes 3 V or more after a lapse of 10.sup.6 seconds of the disturbance time and becomes higher than the Ref voltage of 3 V of the sensing circuit in the reading stage, as a consequence of which data \"0\" is erroneously detected as data \"1\", resulting in erroneous reading.\nA method for alleviating she substrate disturbance as described above is disclosed in the prior art reference of Japanese Patent Laid-Open Publication No. HEI 10-92958 concerning the AND type flash memory. In this specification, a description of the AND type flash memory will be given on condition that memory cell characteristics of this erasing method are similar to the characteristics of the aforementioned ACT type flash memory in order to clarify the problem of the erasing method of the non-volatile semiconductor storage device. That is, the application voltage conditions in the program operation and the erase operation are assumed to be similar to those of Table 1.\nAs shown in FIG. 10, an AND type flash memory has an array structure in which memory cells M00, M01, . . . , M10, M11, . . . are arranged in a matrix form, each of the memory cells being constructed of a floating gate type field-effect transistor capable of electrically writing and erasing information. Word lines WL0 through WL31 and WL32 through WL63 are connected to control gates of the memory cells M00, M01, . . . , M10, M11, arranged in an identical row. The memory cells M00, M01, . . . whose control gates are connected to the word lines WL0 through WL31 belong to a BLOCK0. The memory cells M10, M11, . . . whose control gates are connected to the word lines WL32 through WL63 belong to a BLOCK1. In the memory cells M00, M01, . . . of the BLOCK0, sub-bit lines SBL00 through SBL04094 are jointly connected to drains of the memory cells arranged in an identical column, while source lines SL00 through SL04094 are jointly connected to sources of the memory cells arranged in an identical column. Main bit lines BL0 through BL4094 are connected to the sub-bit lines SBL00 through SBL04094 via selection transistors ST00A through ST04094A, while a selection gate signal line DSG0 is connected to the gates of the selection transistors ST00A through ST04094A. A common source line SL is connected to the source lines SL00 through SL04094 via selection transistors ST00B through ST04094B, while a selection gate signal line SSG0 is connected to the gates of the selection transistors ST00B through ST04094B. In the memory cells M10, M11, . . . of the BLOCK1, sub-bit lines SBL10 through SBL14094 are connected to the drains of the memory cells of an identical column, while source lines SL10 through SL14094 are connected to the sources of the memory cells of an identical column. The main bit lines BL0 through BL4094 are connected to the sub-bit lines SBL10 through SBL14094 via selection transistors ST10A through ST14094A, while a selection gate signal line DSG1 is connected to the gates of the selection transistors ST10A through ST14094A. The common source line SL is connected to the source lines SL10 through SL14094 via selection transistors ST10B through ST14094B, while a selection gate signal line SSG1 is connected to the gates of the selection transistors ST10B through ST14094B.\nIn the AND type flash memory having the aforementioned construction, the case is herein considered where information of memory shells M00, M01, . . . , in the selected block BLOCK0 is subjected to erasing.\nA high positive voltage Vpp (+10 V, for example) is applied to the word lines WL0 through WL31 of the selected block BLOCK0, and a voltage Vnv (-8 V, for example) is applied to all the main bit lines BL0 through BL4094 and a semiconductor substrate (well). A reference voltage Vss (0 V, for example) is applied to the source lines SL00 through SL4094 via the common source line SL. In this stage, the reference voltage Vss (0 V, for example) is applied to the selection gate signal line DSG0 and the voltage Vnv (-8 V, for example) is applied to the selection gate signal line SSG0. Then, the selection transistors ST00A through ST04094A whose gates are connected to the selection gate signal line DSG0 are turned on, so that the voltage Vnv (-8 V, for example) is outputted to the sub-bit lines SBL00 through SBL04094. The selection transistors ST00B through ST04094B whose gates are connected to the selection gate signal line SSG0 are turned off, so that the diffusion source lines SL00 through SL04094 are brought into the floating state. By this operation, the channels of the memory cells M00, M01, . . . of the selected block BLOCK0 are turned on, by which the channel layer comes to have a voltage of -8 V to inject electrons into the floating gate. Consequently, the threshold voltage of the memory cells M00, M01, . . . of the selected block BLOCK0 increases to end the erasing.\nOn the other hand, in the unselected block BLOCK1, the reference voltage Vss (0 V) is applied to the word lines WL32 through WL63 connected to the control gates of the memory cells M10, M11, . . . The voltage Vnv (-8 V) is applied to the selection signal gate line DSG1, so that the selection transistors ST10A through ST14094A whose gates are connected to the selection gate signal line DSG1 are turned off, and consequently the sub-bit lines SBL10 through SBL14094 are brought into the floating state. By applying the voltage Vcc (+3 V) to the selection gate signal line SSG1 and turning on the selection transistors ST10B through ST14096B whose gates are connected to the selection gate signal line SSG1, by which the reference voltage Vss (0 V) is outputted to the source lines SL10 through SL14094 formed from the diffusion layer via the common source line SL. By this operation, a depleted layer is formed instead of a channel layer in the semiconductor substrate (well) just below the tunnel oxide film of the memory cells M10, M11, . . . of the unselected block BLOCK1. For the above reasons, the electric field between the floating gate and the semiconductor substrate (well) is alleviated, by which the substrate disturbance is alleviated.\nHowever, in the aforementioned AND type flash memory, some of the sub-bit lines SBL10 through SBL14094 in the floating state come to immediately have the voltage of -8 V when the voltage of -8 V is applied to the semiconductor substrate (well) due to the diffusion leak (including minute defects) and so on.\nFor example, the case is considered where a leak current of 0.1 .mu.A exists in the sub-bit line formed from the diffusion layer. This is because the threshold voltage of a memory cell is generally defined as the voltage of the word line when the current flowing through the memory cell is 1 .mu.A in the case of the flash memory, and there are practically many sub-bit lines through which the leak current of about 0.1 .mu.A flows. In the case of the flash memory, it is practical that the leak current of the diffusion layer is not so much reduced by comparison with the DRAM.\nIn this case, the voltage of -8 V is applied to the semiconductor substrate (well), and a time Ts during which the sub-bit line that should be in the floating state comes to have the voltage of -8 V is expressed by: ##EQU1##\nwhere C: sub-bit line capacitance (0.02 pF)\nV: sub-bit line voltage (-8 V) PA1 Ir: leak current (0.1 .mu.A) PA1 main bit lines each connected to an associated sub-bit line so as to form a layered structure together with the associated sub-bit line, wherein: in an erase operation of a selected block of the memory cell array, a first negative voltage is applied to the semiconductor substrate, a first positive voltage is applied to the word lines of an unselected block of the memory cell array, and a reference voltage is applied to the sub-bit lines of the unselected block so that memory cells in a low threshold voltage state within the unselected block are turned on, and that a channel layer formed in each of the memory cells which have been turned on comes to have the reference voltage.\nNormally, the erase pulse time is about 1 ms, and therefore, the sub-bit line comes to sufficiently have the voltage of -8 V. In this case, the channel layer is formed in the vicinity of the sub-bit line, as a consequence of which a high electric field is generated between the floating gate and the channel layer (-8 V) in the portion, and electrons are injected into the floating gate, increasing the threshold voltage. In practice, when the memory cell channel layer is sufficiently turned on (i.e., when the channel layer is formed between the source and the drain), the source side of the sub-bit line, which is connected to the common source line SL, comes to have a voltage of 0 V, and therefore, the sub-bit line comes to have a voltage (-6 V, For example) higher than the voltage of -8 V, instead of the voltage of -8 V. However, if the sub-bit line comes to have a voltage higher than -6 V , then the channel layer is cut off by a back gate effect. Therefore, the sub-bit line does not come to have a voltage higher than -6 V (the absolute value is not reduced). The voltage of this sub-bit line differs depending on the threshold voltage and so on of the memory cell.\nTherefore, the nonvolatile semiconductor storage device erasing method described above has an disadvantage that the substrate disturbance cannot be stably alleviated."}
{"text": "1. Field of the Invention\nThe present invention relates to a contents playback apparatus for playing back stored contents based on a desired function, a control method for such a contents playback apparatus, and an electronic device.\n2. Description of the Related Art\nIn recent years, there have been developed small portable contents playback apparatus for storing contents such as music pieces, photographs, and moving images, in a semiconductor memory and playing back the stored contents according to instructions from the user. Particularly, as the storage capacity of semiconductor memories has increased and the cost thereof has decreased, contents playback apparatus that are very small in size and can store a large amount of contents have emerged and are finding a wider range of applications.\nSince the contents playback apparatus can store a large amount of various contents, they incorporate various playback functions to be able to quickly find and play back favorable data of the user. For example, if a contents playback apparatus has read music data from a CD (Compact Disc) or the like and stored the read music data, then the user can select, from the functions of the contents playback apparatus, a function to retrieve desired music pieces sorted by artist or album and play back the retrieved music pieces, a function to narrow down desired music pieces in a certain range such as a play list set by the user and play back the chosen music pieces, or a shuffle playback function to play back music pieces in a sequence which differs each time they are played back.\nHeretofore, a mechanism with a dial-type operation element for easily selecting and deciding on contents to be played back and a playback function in a contents playback apparatus has been disclosed in Japanese Patent Laid-open No. 2003-84902.\nThe mechanism with the dial-type operation element for selecting contents to be played back includes an electronic switch disposed in a casing. The operation of the operation element that is manipulated by the user is mechanically transmitted to the electronic switch. Therefore, a certain link mechanism needs to be provided between the operation member and the electronic switch. However, the link mechanism presents an obstacle to efforts to reduce the size of the casing and also to efforts to effectively utilize the space in the casing. Another problem is that as the link mechanism is provided between the operation member and the electronic switch in the casing, the casing needs to have a hole defined in a wall thereof for passage of the link mechanism therethrough, and hence cannot be made sufficiently water-resistant."}
{"text": "Electronic devices and systems can include various sub-elements that can receive individualized input power. This input power can be applied or removed as-needed, such as to perform within low power architectures that place various components into low-power or off modes when not in use. Individual portions of microprocessors or system-on-a-chip (SoC) devices can also include separate power domains that can be powered on and off independently of each other. These techniques can be referred to as power gating, and are often employed to conserve power in electronic devices, such as in computers, handheld devices, smartphones, gaming systems, and the like."}
{"text": "Reflectors for use with LED light sources typically are constructed from conductive, reflective materials, such as aluminum or vacuum metalized substrates. A number of disadvantages exist when using reflectors of this type. For instance, the use of conductive materials in the entire reflector generally requires that an isolation gap be maintained between the reflector and LED light source. The isolation gap required is based on a minimum creepage distance to protect against electric discharges on or close to an insulation surface and a minimum clearance distance to prevent dielectric breakdown between conductive parts by the ionization of air. This requirement for the isolation gap results in a reflector that is too far from the LED light source. The resultant gap reduces the ability to control light being emitted from the light source as efficiently and effectively, as some light is typically lost along the gap. In addition, in instances where the reflector needs to be easily and quickly replaced, the coaxial orientation and position of the reflector must be maintained after the reflector is replaced so that the beam control and light distribution is not affected.\nIn the case of metalized reflectors, these reflectors can include a plastic piece that is injection molded, and then metalized with a conductive material to achieve a reflective surface. A coating, such as a lacquer coating, must be applied to the metalized surface thereafter to protect the metallization. However, the coating generally degrades over time and the reflectivity diminishes as a result. In general, as the coating degrades, the color accuracy and total system efficiency is impacted. In addition, these metalized reflectors are conductive."}
{"text": "Papaya (Carica papaya L.) is an important tropical fruit crop with a yielding potential of approximately 45 tons/hectare, which is normally consumed fresh and is valued as a health food because it is rich in vitamins C and A. Papaya is widely grown in Brazil, Australia, South Africa, South-East Asia, Hawaii, India and other tropical areas. Papaya is a polygamous species, and sex inheritance is controlled by a single locus with multiple alleles. There are three sex forms of papaya trees, i.e. hermaphrodite or bisexual, pistillate or female, and staminate or male. The sex types of papaya may not be identified according to the phenotype of the juvenile plant or other chemical or biochemical methods. The male, female and hermaphroditic flowers of papaya are distributed on separate papaya plants and sex types are revealed only after flowering. Methods to identify sex type at juvenile stage have been studied (Bojappa and Singh 1974, Choudhri et al. 1957, Parasnis et al. 1999, 2000; Singh et al. 1977, Sondur et al. 1996, Somsri et al. 1998). Storey proposed that the sex of papaya is determined by three homologous gene complexes on sex chromosome (W. B. Storey, J. Hered. 44, 70–78, 1953; and W. B. Storey, Crop Plants, 21–24. Wisley, N.Y.). The genes are so tightly linked that no crossing over occurs among them; thus the complexes are transmitted to offspring as if they are single gene alleles with pleiotrophic effects on phenotypic expression. The genotypes of the male, hermaphroditic and female plants are M1m, M2m, and mm, respectively. Genotypes with homozygous dominant alleles are lethal (W. B. Storey, J. Hered. 44, 70–78, 1953).\nThe hermaphroditic papaya, bearing perfect flowers and producing fruits shaped from long-cylindrical to ellipsoidal, is preferred by the markets in Hawaii, Japan, South-East Asia and Taiwan. In addition, the consumers and farmers prefer hermaphrodite papaya because it has small seed cavity and is easier to package. The cross between two hermaphroditic papayas will yield a ratio of 2:1 hermaphrodite to female papaya. Therefore, the papaya growers usually plant at least two seedlings in each hole in the field, and later remove the females at flowering, a practice that is time-consuming and wasteful. Therefore, there is a need to develop a technology of obtaining all hermaphrodite papaya."}
{"text": "The present invention relates to devices for selecting numbers or combinations of numbers in games of chance.\nGames of chance, wherein the players choose their own number or combinations of numbers, have become increasingly popular over the past few years. Governments, in order to increase revenue, frequently operate this type of game or lottery. The player receives a ticket with his selected number combination and the operator chooses a winning number on a weekly, daily or other basis. If the player's selected number combination matches the winning number, the player is a winner.\nSome prizes in State lotteries have reached astronomical proportions. For example, one lottery winner was recently awarded over 16 million dollars. To cash in on this opportunity, players often select several combinations for one drawing. To do this, unprepared players often have to make a large number of selections in a relatively short time--that time being when the operator says, \"Next! And what is your number?\"\nPlayers typically use various mathematical formulas or special numbers associated with birthdays, anniversaries and other meaningful dates. However, sometimes players lose faith in their numbers and look for any crutch or sign to assist them in winning.\nSome devices have been previously developed which would allow random selection of combinations of numbers. Most of these selectors are hand-held and most are hand-agitated. Typically, an enclosed reservoir or chamber of number indicators, usually small numbered balls, is agitated or shaken by the player to insure a random array of the number indicators. Then the player opens a gate or slide which permits the number indicators to leave the holding reservoir or chamber to be physically withdrawn from the device or to be viewed within a display section of the device.\nWhile such devices can select random number combinations relatively easily, none have gained widespread acceptance. One reason is that, although the devices are easily portable, lottery players rarely bring an apparatus to the site of the number combination selection. Additionally, these prior devices would be impractical if they were offered at the lottery \"store\" because their hand operation renders them subject to rapid wear or abuse, especially when their numbered balls have to be plucked and held by a user.\nAccordingly, it is the primary object of the present invention to provide an improved number-combination selector which randomly selects number combinations from a closed reservoir or chamber without having to be hand agitated to produce the combinations.\nIt is another object to provide a lightweight, portable number-combination selector, wherein the number indicators are selected and stored for display in a clear sealed chute and then returned to the storage or agitation chamber without being touched by the player.\nIt is yet another object to provide a number-combination selector which is commensurate with the above-listed objects and comprises an adaptor hood that fits over a standard, underlying bingo-ball blower.\nIt is a still further object to provide a number-combination selector which is easily amenable to being operated as a commercial vending machine, wherein the player inserts a coin to start the device for a specific time period during which the player selects the combinations.\nThe above and other objects and advantages of this invention will become more readily apparent when the following description is read in conjunction with the accompanying drawings."}
{"text": "The present invention relates to a display device having a scale.\nDisplay devices in round shape are already known in which a rotating pointer marks the values to be indicated\nFurthermore, digital display devices are known in which the value to be indicated is represented by a sequence of digits.\nIt has been found that an indication in the form of an analog pointer means has the advantage over said digital means that it can be immediately noted visually by the user. Analog pointer means in the form of a rotating pointer have, however, the disadvantage that a substantially round dial must be selected for the direct visual perception of the value indicated.\nThe object of the present invention is, therefore, to provide a display device which has the advantages of an analog visual indication, with which the value indicated can be read conveniently even in the case of scales having smaller dimensions."}
{"text": "1. Field of the Invention\nThis invention relates to the air slide delivery of powders of all types and to the transport of particulate matter by a gas.\n2. Description of the Prior Art\nAir slide systems with filters to remove transport air from the vessel to which the particulate matter is delivered are old. These systems often require costly \"bag house\" filters to handle the large volumes of transport air involved."}
{"text": "1. Field of the Invention\nThe invention relates to a diffractive optical element having a multiplicity of binary blazed diffraction structures. The diffractive optical element is particularly intended for use in microlithographic projection exposure apparatus.\n2. Description of the Prior Art\nConventional blazed gratings have diffraction structures of triangular, in particular sawtoothed cross section which extend mutually parallel with a spacing equal to the grating constant g. One edge of the diffraction structures, the blaze edge, has an inclination with respect to the base surface of the grating such that the reflection or refraction law is satisfied for one diffraction order of the incident light, and the majority of the intensity of the diffracted light is therefore contained in the order favoured by the blaze edge. The traditional method of producing such blazed gratings consisted in scratching the diffraction structures in a master grating with the aid of diamonds and making corresponding copies of this master grating. This mechanical method is highly elaborate, on the one hand, and on the other hand it encounters limitations with very short wavelengths of the light for which the grating is intended to be used, since the structures to be produced are too small.\nEfforts have therefore been made to employ the process technology used for the production of semiconductor components, in which a substrate is coated with photoresist, exposed, subsequently developed and etched, in order to produce the diffraction structures of blazed gratings. The approach firstly involved using successions of such process cycles to achieve diffraction structures which are supposed to approximate the blaze edge by a stepped edge. If four such steps are used, for example, then diffraction efficiencies of more than 80% can be achieved in the first order. With a further process cycle, eight stages are obtained by which a first-order diffraction efficiency of about 95% can be achieved. In general, 2n steps can be produced by using n process cycles. With increasing n, the stepped profile of the edge becomes closer and closer to the sawtooth profile of ideal blazed gratings in conventional, mechanically produced gratings, the diffraction efficiency of which is 100% in the first order according to scalar theory. The production of such a grating, however, is cost-intensive and error-prone because it is necessary to carry out the process cycle repeatedly.\nAttempts have also been undertaken to simulate the blaze profile of the diffraction structures by using binary structures whose dimensions are smaller than the wavelength of the electromagnetic radiation for which the grating was defined. These attempts are based on the fact that light is no longer diffracted at the small substructures, but can only be scattered. This leaves only the zeroth diffraction order which picks up the effect of the substructures merely in the form of a local effective refractive index in phase gratings, or merely in the form of a local shade of grey in amplitude gratings.\nA first example of such a binary blazed grating is described in the article by Joseph N. Mait et al. “Diffractive lens fabricated with binary features less than 60 nm”, Optics Letters, 15 Mar. 2000, pages 381 et seqq. The substructure used here is a multiplicity of lines, all of which extend parallel to the diffraction structure and whose spacing is less than the effective wavelength.\nThe article by Philippe Lalanne et al. “Design and fabrication of blazed binary diffractive elements with sampling periods smaller than the structural cut off”, J. Opt. Soc. Am. A, May 1999, pages 1143 et seqq. describes blazed diffractive elements of the type mentioned in the introduction, in which the diffraction structures are resolved into individual substructures consisting of rectangular or square pillars. Different “fill factors” can be achieved by varying the pillar width for a predetermined pillar spacing, and this corresponds to a local variation of the effective refractive index. As an alternative, the pillars may also be arranged at different spacings with a constant width.\nA common feature of all these attempts to produce binary blazed diffractive optical elements is that the substructures are minutely configured and have a very high aspect ratio (structure height to structure width). They are therefore technologically highly elaborate and expensive to produce, and cannot be made with sufficient accuracy."}
{"text": "The present invention relates to fuel vapor vent valves employed in fuel tanks filled with highly volatile fluids such as gasoline or mixtures of gasoline and alcohol and particularly relates to valves employed in motor vehicle fuel tanks.\nCurrently passenger cars and light truck vehicles employ fuel vapor storage devices connected to the vehicle fuel tank through a float operated valve which controls the flow of vapor in the dome above the liquid fuel level to a storage device. Currently, such valves are required to prevent the escape of liquid fuel in the event of overfilling the tank or angular displacement of the vehicle including rollover conditions.\nWith the advent of molded plastic fuel tanks, it has been found difficult to provide an economical design for the vent valve and the attachment of the valve to the tank. Heretofore, such vent valves had been installed through an access opening in the tank which requires sealing in a manner sufficient to prevent escape of vapor and permeation of the vapor through the material of the vent valve and the tank wall.\nIt has been found that the material required to withstand continuous exposure to the liquid fuel and vapor has the propensity to be permeable to the fuel vapor. This problem has been addressed by molding the tank wall of layers of different materials with a vapor impervious barrier layer embedded in the material of the tank wall.\nIf an access opening is formed in the top of the tank for installation of a vapor vent valve, the vapor barrier continuity is broken and the potential for localized vapor permeation has resulted.\nAccordingly, it has been proposed to install the vapor vent valve on the interior of the fuel tank without forming an access opening in the upper wall of the tank.\nThe aforesaid proposal of mounting a vapor vent valve on the interior of the tank has been complicated by the use of high density polyethylene (HDPE) material for molding of the fuel tank. For economical installation of the vent valve on the inside of the fuel tank, the use of HDPE material for the valve has been required in order to permit securing the valve to the wall of the tank by weldment. However, HDPE material has been found not satisfactory for the structural components of the float operated valve and this has resulted in difficulties in designing and manufacturing a valve for interior installation in the tank.\nThe present invention provides a unique and novel technique for installing a float operated fuel vapor vent valve inside a fuel tank and securing the valve to the tank wall by weldment. The present invention provides for a cup-shaped or U-shaped attachment member with the valve received therein and recessed below the rim of the attachment member which is secured to the inner surface of the tank wall by weldment. The attachment member secures and retains the valve in the desired position and orientation in the tank. The arrangement of the present invention thus permits the body structure of the valve to be formed of a desired material different from the material of the attachment member which is required to be the same as the tank wall material in order to facilitate attachment by weldment."}
{"text": "1. Field of the Invention\nThe present invention relates to a nonvolatile semiconductor memory device having a memory cell array in which memory cells each comprising a variable resistance element storing information by the change of an electric resistance are arranged in a row direction and a column direction, and more particularly, to a technique for preventing and suppressing the deterioration of stored data due to the reading operation of the memory cell array.\n2. Description of the Related Art\nRecently, there has been proposed a variable resistance type of memory element (referred to as the variable resistance element hereinafter) having a two-terminal structure in which a metal oxide is sandwiched by conductors serving as electrodes, and capable of changing its electric resistance reversibly by applying a voltage pulse. Various kinds of variable resistance elements are proposed and disclosed by combining oxide materials and electrode materials (or example, refer to document 1: Japanese Laid-Open Patent Publication No. 2004-087069, document 2: Liu, S. Q. et al., “Electric-pulse-induced reversible Resistance change effect in magnetoresistive films”, Applied Physics Letter, Vol. 76, 2749, in 2000, document 3: Seo, S. et al., “Reproducible Resistance Switching in polycrystalline NiO films”, Applied Physics Letters, Vol. 85, 5655, in 2004, document 4: Sim, H et al., “Resistance-switching characteristics of polycrystalline Nb2O5 for nonvolatile memory application”, IEEE Electron Device letters, Vol. 26, 292, in 2005, document 5: Sawa, A. et al., “Hysteretic current-voltage characteristics and resistance switching at rectifying Ti/Pr0.7Ca0.3MnO3 interface”, Applied Physics Letters, Vol. 85, 4073, in 2004, and document 6: Fujii, T. et al., “Hysteretic current-voltage characteristics and resistance switching at an epitaxial oxide Schottky junction SrRuO3/SrTi0.99Nb0.01O3”, Applied Physics Letters, Vol. 86, 12107, in 2005), and each provides distinctive electric characteristics and varies in operation mechanism. Every variable resistance element uses a reversible resistance changing operation (referred to as the “switching operation” occasionally hereinafter) and can be used as a new nonvolatile semiconductor memory device by relating information to a resistance value and reading the resistance value or a current corresponding the resistance value. Here, the information includes binary digital data, a multilevel digital data, analog data and the like, and the high resistance state and the low resistance state are stored as the binary digital data “1” and “0”, the multilevel digital data can be stored using a middle resistance value between the high resistance state and the low resistance state, or the analog data may be stored.\nThere can be constituted a nonvolatile semiconductor memory device by forming a memory cell array in which memory cells comprising the variable resistance element and storing information by the change of the electric resistance of the variable resistance element are arranged in a row direction and column direction in a matrix state, and providing a circuit for controlling programming, erasing and reading operation of data for each memory cell of the memory cell array in the vicinity of the memory cell array.\nThe constitution of the memory cell comprising the variable resistance element includes a case where each memory cell comprises a series circuit consisting of the variable resistance element and a transistor as a cell-access element (1T/1R-type memory cell), a case where each memory cell comprises a series circuit consisting of the variable resistance element and a diode as a cell-access element (1D/1R-type memory cell), a case where each memory cell comprises a variable resistance element only (1R-type memory cell) and the like. The 1T/1R-type memory cell and its memory cell array, for example are disclosed in the document 1 by the applicant of this application (refer to FIG. 1, for example). The 1D/1R-type memory cell, for example is disclosed in Japanese Laid-Open Patent Publication No. 2004-260162 by the applicant of this application (refer to FIG. 1, for example). The 1R-type memory cell, for example is disclosed in Japanese Laid-Open Patent Publication No. 2005-32401 (refer to FIG. 4, for example).\nWhen data is read from the memory cell comprising the variable resistance element, a bias voltage is applied to the variable resistance element to flow a reading current and the resistance value of the variable resistance element is determined by the amount of the current, so that the data is read. Therefore, regardless of the constitution of the memory cell, a predetermined bias voltage is applied to the variable resistance element in the reading operation. When a phenomenon in which the resistance value of the variable resistance element is changed a little by the bias voltage applied at the time of this reading operation (referred to as the “reading disturbance” occasionally hereinafter) is repeated, recorded information could be lost in the worst case. Therefore, it is necessary to reduce the degree and frequency of the reading disturbance as much as possible.\nAs described above, although there are various kinds of nonvolatile variable resistance element capable of changing the electric resistance reversibly by applying the voltage pulse, the behavior of the reading disturbance in the variable resistance element is not clear.\nThe inventors have found that in the case where the variable resistance element showing the rectifying characteristics disclosed in the document 5 (Sawa, A. et al.) or the document 6 (Fujii, T. et al.) is used, when a reading voltage whose absolute value is not more than a programming voltage is applied to the variable resistance element continuously, the resistance value of the variable resistance element is changed and the resistance value is considerably changed depending on the polarity of the reading voltage. In addition, the variable resistance element showing the rectifying characteristics denotes that the variable resistance element itself has the rectifying characteristics and does not mean that when the memory cell comprises a series circuit consisting of the variable resistance element and a diode as a cell-access element, the memory cell has the rectifying characteristics.\nFIG. 1 shows current-voltage characteristics in a high resistance state and a low resistance state of the variable resistance element disclosed in the document 5 (Sawa, A. et al.) and comprising three layers Ti/Pr0.7Ca0.3MnO3 (PCMO)/SrRuO3 (SRO) manufactured by a similar method to that disclosed in the document 5 (Sawa, A. et al.). An upper electrode is Ti and an applied voltage in FIG. 1 is the potential of the upper electrode based on a lower electrode. Referring to FIG. 1, since a negative bias current when a negative voltage is applied (at the time of negative bias) is larger than a positive bias current when a positive voltage is applied (at the time of positive bias), forward bias is provided for the rectifying characteristics at the time of negative bias while reverse bias is provided for the rectifying characteristics at the time of positive bias. In addition, when the potential of the lower electrode based on the upper electrode is defined as the applied voltage, the above relation is reversed. Furthermore, the forward bias is defined by an applied voltage polarity in which a larger current flows to the variable resistance element.\nIn addition, according to the current-voltage characteristics shown in FIG. 1, when the current-voltage characteristics in the high resistance state is compared with that in the low resistance state, a current difference is largely provided in both forward bias and reverse bias, so that the high resistance state and the low resistance state can be determined in the reading operation in both forward bias and reverse bias.\nHowever, the inventors of the present invention has found that the degree of the reading disturbance is considerably different between the reading operation in the forward bias (forward reading) and the reading operation in the reverse bias (reverse reading). FIG. 2 shows graphs in which the change in resistance value is plotted with a reading voltage applying time (reading voltage pulse applying number of times) when the forward reading and reverse reading are performed for variable resistance elements in the low resistance state and in the high resistance state. The change in resistance value is shown relatively assuming that the resistance value just after the variable resistance element becomes the low resistance state or high resistance state is set to 1, which denotes that the characteristics becomes undesirable as the resistance value (relative value) becomes far from 1. It can be seen from FIG. 2 that the resistance value change is larger when the reverse reading is performed in the low resistance state than the other case. In addition, since the resistance value in the above reading operation tends to increase, when the same reading operation is continued, the resistance state is changed from the low resistance state to the high resistance state, so that recorded information is lost.\nThus, it is clear from the above experimental result that the reading disturbance phenomenon is such that the data stored in the memory cell, that is, the resistance value is changed with the voltage applying time (number of times for applying a pulse) in the reading operation. Especially, the resistance value of the variable resistance element is considerably changed when the reading operation is performed by applying the reading voltage to the variable resistance element in the low resistance state in the reverse bias, so that when the same reading operation to the same memory cell is repeated, stored data could be completely lost and could not read in the worse case.\nFurthermore, since in the case of the memory cell array comprising the 1R-type memory cell, the reading voltage is also applied to the selected memory cell that is not to be read but shares the word line or bit line with the memory cell to be read, the above reading disturbance phenomenon appears more notably, so that it is highly necessary to prevent the reading disturbance phenomenon as compared with the other memory cell types."}
{"text": "Such a mounting device is known from U.S. Pat. No. 3,579,840. Therein, in each case a fixed clamping jaw as well as an adjustable clamping jaw on the opposite side are arranged on a ring-shaped holding device. On the bottom side of the ring-shaped holding device, between the two clamping jaws, a bar-shaped recoil cleat is provided for the engagement in a corresponding continuous transverse groove on a base portion which can be attached to the handgun. By means of the recoil cleat, the forces acting at the time of the firing are absorbed, and as a result an improved precision can be achieved. However, as a result of the continuous transverse grooves, the attachment device is also weakened, and thus the stiffness is affected."}
{"text": "The present invention relates to a thyristor, and more particularly, to a mechanism for accurately positioning its light-sensitive area and trigger guide in registry.\nThis mechanism is hereunder described with particular reference being made to a light pulse triggered thyristor. A typical example of a conventional light pulse triggered thyristor is shown in cross-section in FIG. 1, wherein a thyristor element 1 has a trigger section or light-receiving area 1ain the center of its upper surface. A reinforcing disk 2 made of a metallic material having a coefficient of thermal expansion close to that of the element 1 is secured to the lower surface thereof. A first electrode unit 3 is pressed against an electrode that is formed to surround the light-receiving area 1a on the upper surface of the element 1. A second electrode unit 4 is pressed against the lower surface of the reinforcing disk 2. An insulating tube 5 for positioning and reinforcing the metal disk 2 is made of a material like alumina ceramic and peripherally surrounds the element 1, the disk 2, and the electrode units 3 and 4. The inside diameter of the insulating tube 5 is slightly larger than the outside diameter of the disk 2. Upper and lower annular metal plate flanges 6 and 7, respectively, are hermetically secured to the upper and lower sides of the insulating tube 5. The inner peripheral surfaces of the flanges 6 and 7 are similarly secured to the outer peripheral surfaces of the electrode units 3 and 4, respectively. A transverse groove 8 is provided in the first electrode unit 3 on the side thereof facing the element 1. Said groove 8 extends radially inward from the outer peripheral surface of the first electrode unit 3 to a point beyond the light-receiving area 1a and is wider in a direction perpendicular to the section of FIG. 1 than the diameter of the light-receiving area 1a. A transverse through-hole 9 is made in an area of the insulating tube 5 corresponding to the groove 8. A light guide 10 for directing external trigger light signals to the light-receiving area 1a comprises a glass rod permitting a high degree of light transmission. The outer peripheral surface of the portion of said guide 10 which is closer to its outside end is sealed to the inner wall of the through-hole 9 so that the inside end of said guide is positioned along the central axis of the insulating tube 5 at a point close to the light-receiving area 1a. Said outside end passes through the transverse groove 8 and the through-hole 9 to extend outside of the insulating tube 5.\nIn the apparatus of FIG. 1, the outer peripheral surface of the portion of the light guide 10 which is closer to its outside end is sealed to the inner wall of the through-hole 9 in the insulating tube 5 so that the inside end of the guide 10 is positioned along the central axis of the insulating tube 5. The reinforcing metal disk 2 is held in position by being pressed into the insulating tube 5. Therefore, in order for both the inside end of the light guide 10 and the light-receiving area 1a of the element 1 that is secured to the upper surface of the metal disk 2 to be positioned along the central axis of the insulating tube 5, the element 1 must be fixed to the upper surface of the metal disk 2 so that the center of the light-receiving area 1a is positioned on the line that is normal to the center of said upper surface of the metal disk 2. But this is very difficult to achieve and it sometimes occurs that the center of the light-receiving area 1a is not in registry with the line that is normal to the center of the upper surface of the metal disk 2. If this misalignment occurs, the inside end of the light guide 10 becomes offset with respect to the center of the light-receiving area 1a and the amount of light being transmitted from the inside end of the light guide 10 toward the light-receiving area 1a is not sufficient to retain the desired high sensitivity of triggering by light pulses."}
{"text": "The present invention relates generally to gene expression and specifically to a novel enhancer element that increases the rate of transcription of a gene operably linked thereto, particularly in plants.\nGenes are regulated in an inducible, cell type-specific or constitutive manner. There are different types of structural elements which are involved in the regulation of gene expression. Cis-acting elements, located in the proximity of, or within genes, serve to bind sequence-specific DNA binding proteins, i.e., trans-acting factors. The binding of proteins to DNA is responsible for the initiation, maintenance, or down-regulation of gene transcription.\nCis-acting elements which control genes include promoters, enhancers and silencers. Promoters are positioned next to the transcription start site and function in an orientation-dependent manner, while enhancer and silencer elements, which modulate the activity of promoters, may be flexible with respect to their orientation and distance from the transcription start site.\nAn example of a specifically regulated gene in plants is phenylalanine ammonia-lyase (PAL), which catalyzes the deamination of phenylalanine to cinnamic acid, the precursor of a wide variety of natural products based on the phenylpropane skeleton. During vascular development, PAL is selectively expressed in differentiating xylem cells associated with deposition of the structural polymer lignin. Lignin, the second most abundant biopolymer after cellulose, is the major structural cell wall component of cells forming vessels in plant tissue (xylem). The xylem is responsible for movement of water and inorganic solutes from plant roots to plant shoots. PAL, genes are expressed at correspondingly high levels in differentiating xylem.\nThe ability to artificially regulate the rate of gene expression provides a means of producing plants with new characteristics. There are numerous situations in which increased levels of gene expression, including increased endogenous gene expression, may be desirable. Such situations include, for example, production of protein plant products for agricultural or commercial purposes.\nThe present invention provides a novel repeat element which functions as a non-specific enhancer. In other words, the invention enhancer element does not affect the intrinsic specificity of a promoter associated with the enhancer element. Instead, the enhancer element boosts the activity of the promoter thereby resulting in a desired level of expression of a gene associated with the promoter. A novel transcription factor, palindromic element binding factor (PABF), which binds to the novel repeat element is also provided.\nIn a first embodiment, the invention provides an enhancer element comprising an isolated nucleotide sequence consisting of at least the sequence (AATT)n, where n=2, and preferably from about 2 to about 20. The sequence (AATT)n has cis-acting, non-specific, enhancer activity. In one aspect, the invention provides a method for increasing expression of a gene in a cell comprising operably linking a (AATT)n repeat element to a heterologous promoter which is operably linked with the gene, thereby permitting increased expression of the gene.\nIn another embodiment, the invention provides a substantially purified palindromic element binding factor (PABF) polypeptide characterized as having a molecular weight of approximately 67 kDa, as determined by SDS-PAGE, binding to a (AATT)n repeat element, where n=2, and having a H1 histone domain, a glutamine rich domain and a high mobility group (HMG) I/Y domain. PABF acts as a transcription factor and binds to the (AATT) repeat element of the invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a connecting structure for a striking plate of a golf club head. In particular, the present invention relates to a connecting structure for connecting a striking plate to a body of a golf club head for simplifying assembling and positioning for a subsequent welding procedure.\n2. Description of Related Art\nTaiwan Patent Publication No. 327606 discloses a method for connecting a golf club head body and a striking plate, both made of metal. The golf club head body includes a recess for engaging with the striking plate. The recess includes a shoulder on which a welding material is placed. The striking plate is inserted into the recess and presses against the shoulder to cause deformation of a protruded portion on an inner edge of the shoulder, thereby filling the welding material into a gap defined between the striking plate, the shoulder, and the protruded portion. The welding material is in the form of metal powder and has a melting point lower than that of the golf club head body and that of the striking plate. The combination of the golf club head body and the striking plate is placed into a vacuum furnace or an inert gas atmosphere in a high temperature furnace and then heated at a temperature higher than the melting point of the welding material and lower than the melting point of the golf club head body and lower than the melting point of the striking plate. The molten welding material fills the tiny gaps between the golf club head and the striking plate by capillary action. A solid golf club head without welding marks on appearance is obtained after cooling.\nNevertheless, the gap between the striking plate, the shoulder, and the protruded portion can receive a limited amount of welding material, with a portion of the welding material filling the tiny gaps between an inner perimeter delimiting the recess and an outer perimeter of the striking plate. Thus, the remaining welding material is insufficient to fill the gap between the striking plate, the shoulder, and the protruded portion. Cavities are formed accordingly. As a result, when the striking plate is subjected to a striking stress and thus elastically deforms, cracks are apt to be generated in the welding areas. Further, the gap between the outer perimeter of the striking plate and the inner perimeter delimiting the recess must be precisely controlled to assure the welding material for braze welding to fill the tiny gaps by capillary action. Hence, additional equipment is required for milling the golf club head body and the striking plate so as to precisely control the tolerance of the gap regardless of the process for manufacturing the golf club head body and the striking plate (such as precision casting). The overall time for manufacturing the golf club head and the manufacturing cost are both increased, which is detrimental to mass production.\nTaiwan Patent Publication No. 469144, a patent of addition of Taiwan Publication No. 327606, proposes an ordinary welding (such as argon welding) along a seam between the striking plate and the golf club head body after braze welding. Finally, the outer surface of the golf club head is subjected to grinding and surface finishing to guarantee the bonding strength by external welding. Nevertheless, the tolerance between the golf club head body and the striking plate requires precise control. Further, the overall time for manufacturing the golf club head and the manufacturing cost are both increased, which is detrimental to mass production."}
{"text": "1. Technical Field\nThe present disclosure relates to printing methods and systems, and more particularly, to a printing method and system to print files from an electronic book (e-book).\n2. Description of Related Art\nNowadays, mobile electronic devices generally can store various files, such as pictures and text files. In order to print a document, a printer driver must be installed on a mobile electronic device. The printer driver corresponds to a first printer. However, if a user wants to print to a second printer, then a different printer driver corresponding to the second printer must be installed. This process must be repeated for the second printer, which is time-consuming and inconvenient."}
{"text": "Many attempts have been made to obtain improved plants for cultivation through breeding programs. A conventional plant breeding program requires as much as ten years to develop a new variety. In addition to the initial hybridization step, several years are typically spent replanting successive generations in order to obtain homozygous plants. An alternative to a conventional plant breeding program is anther culturing in which anthers from one plant are used to pollinate the ovaries of another plant. However, many traits cannot be successfully introduced via such hybridization techniques since the genes for such traits are not found in breeds available for hybridization.\nAn alternative to hybridization is somaclonal variation. This technique involves the use of vegetative plant parts, such as callus tissue, as explant material. For example, callus that develops from vegetative explants of rice frequently regenerates plants which have genetic characteristics not found in the variety from which the explant was originally obtained. These somatic mutants occur at high frequencies, and the percentage of regenerated plants which differ from the starting variety exceeds, for example, 75 percent in rice. This technique is therefore useful for producing genetic variability. Again, however, there are limits to the extent of variation which can be obtained.\nDevelopment of plant genetic engineering began in the early 1940s when experiments were being carried out to determine the biological principle causing formation of crown gall tumors. The tumor-inducing principle was shown to be a bacterial plasmid from the infective organism Agrobacterium tumefaciens. This plasmid has since been characterized in much detail utilizing the currently available techniques of recombinant DNA technology. The bacterium elicits its response by inserting a small fragment of bacterial plasmid into the plant nucleus where it becomes incorporated and functions as a plant gene. This discovery opened the door to using Agrobacterium and its plasmids as vehicles to carry foreign DNA to the plant nucleus. There are, however, limitations to the application of this technique which include: (1) susceptibility to infection with the Agrobacterium plasmid and (2) available tissue culture technology for regeneration of the transformed plants. Thus, there are no successful reports on genetic engineering of monocots such as cereals with Agrobacterium plasmid vectors because of the general inability of Agrobacterium to infect monocots.\nMore recently, other techniques have been used to genetically transform monocots. For example, electroporation of protoplasts of rice, wheat and sorghum to obtain expression of a foreign gene was reported in Ou-Lee et al, Botany, vol. 83, pp. 6815-6819 (1986). Similar plant protoplast electroporation and electroinjection through cell walls and membranes have also been reported for other monocots, and dicots as well. See, Fromm et al, Proc. Natl. Acad Sci USA, vol. 82, pp. 5824-5828 (1985); Hibi et al, J. Gen. Virol., vol. 67, pp. 2037-2042 (1986); Langridge et al, Plant Cell Reports, vol. 4, pp. 355-359 (1985); Fromm et al, Nature, vol. 319, pp 791-793; Shillito et al, Bio/Technology, vol. 3, pp. 1099-1103 (1985); and Okada et al, Plant Cell Physiol., vol. 27, pp. 619-626 (1986). Similarly, direct and chemical-induced introduction of DNA into monocot and dicot cells has been disclosed. See, Lorz et al, Mol. Gen. Genet., vol. 199, pp. 178-182 (1985); Potrykus et al, Mol. Gen. Genet., vol. 199, pp 183-188 (1986); Uchimaya et al, Mol. Gen. Genet., vol. 204, pp 204-207 (1986); Freeman et al, Plant Cell Physiol., vol. 25, no. 8, pp. 1353-1365 (1984); and Krens et al, Nature, vol. 296, pp. 72-74 (1986). Another technique of interest is the injection of DNA into young floral tillers of rye plants reported in de la Pena et al, Nature, vol. 325, pp. 274-276 (1987).\nThe agricultural production of major crops has long been significantly affected by insects and plant pathogens. For example, blight and blast are major diseases of rice plants which can decimate a crop. Some plants cannot be cultivated in certain parts of the world because of the presence of diseases in such locations to which the plants are susceptible. For example, the main diseases in potato are bacterial soft rot and bacterial wilt caused by Erwinia carotovora and Pseudomonas solanacearum, respectively. These diseases are primarily responsible for limiting the growth of potatoes in many areas of Asia, Africa, South and Central America. Moreover, pesticides are becoming increasingly difficult to use in an effective, and yet environmentally acceptable manner. Therefore, it would be desirable to have available for cultivation plants which are resistant to insects and other pathogens.\nIt is well known that the pupae of Hyalophora (a type of silk moth) respond to bacterial infection by the synthesis of mRNAs which culminate in the production of about 15 to 20 new proteins. Lysozyme, the antibacterial protein found in egg white and human tears, and two other classes of antibacterial peptides, called cecropins and attacins, have been purified from Hyalophora humor. These proteins have a rather broad spectrum of activity in that they are effective on many different types of bacteria. Thus, the insects have evolved a rather successful and novel means to fight bacterial infections. Although a traditional immunologist would think this system lacks specificity, the insect has a rather potent arsenal of at least three different antibacterial proteins which may work in different ways to destroy bacterial pathogens. Thus, the invading bacteria is presented with a formidable challenge which is very difficult to circumvent. While a bacterial pathogen may be naturally resistant to one, it is highly improbable that it would be resistant to all three toxins. Although the exact mode of action of the protein toxins is not fully understood, they are generally procaryote specific and appear to be benign to eucaryotic insect cells.\nAs mentioned above, the property of certain peptides to induce lysis of procaryotic microorganisms such as bacteria is well known. For example, U.S. Pat. Nos. 4,355,104 and 4,520,016 to Hultmark et al describe the bacteriolytic properties of some cecropins against Gram-negative bacteria. Quite interestingly, the cecropins described in the Hultmark et al patents were not universally effective against all Gram-negative bacteria. For example, the cecropins described therein lysed Serratia marcescens D61108, but not Serratia marcescens D611. Moreover, cecropins have generally been reported to have no lytic activity towards eucaryotic cells such as insect cells, liver cells and sheep erythrocytes, as reported in the Hultmark patents; International Patent Publication WO/8604356; Andreu et al, Biochemistry, vol. 24, pp. 1683-88 (1985); Boman et al, Developmental and Comparative Immunology, vol. 9, pp. 551-558 (1985); and Steiner et al, Nature, vol. 292, pp. 246-248 (1981).\nOther lytic peptides heretofore known include, for example, the sarcotoxins and lepidopterans. Such peptides generally occur naturally in the immune system of Sarcophaga peregrina and the silkworm, lepidopteran, respectively, as reported in Nakajima et al, The Journal of Biological Chemistry, vol. 262, pp. 1665-1669 (1987) and Nakai et al, Chem. Abst. 106:214351w (1987).\nA number of the antibacterial polypeptides have been found to be useful when the genes encoding therefor are incorporated into various animals. Such transformation of animals with genes encoding therefor are described in U.S. patent application Ser. No. 069,653, filed Jul. 25, 1986, now abandoned by Jesse M. Jaynes, Frederick M. Enright and Kenneth F. White, which is hereby incorporated herein by reference.\nPolynucleotide molecules expressible in E. coli and having the sequence araB promoter operably linked to a gene which is heterologous to such host are also known. The heterologous gene codes for a biologically active polypeptide. A genetic construct of a first genetic sequence coding for cecropin operably linked to a second genetic sequence coding for a polypeptide which is capable of suppressing the biological effect of the resulting fusion protein towards an otherwise cecropin-sensitive bacterium is also described in International Publication WO86/04356, Jul. 31, 1986.\nThe Hultmark et al patents mentioned above also mention that there are no known antibodies to cecropin, indicating a wide acceptability for human and veterinary applications, including one apparently useful application for surface infections because of the high activity against Pseudomonas. Similarly, EPO publication 182,278 (1986) mentions that sarcotoxins may be expected to be effective in pharmaceutical preparations and as foodstuff additives, and that antibacterial activity of sarcotoxin can be recognized in the presence of serum Shiba, Chem. Abstr. 104:230430K (1985) also mentions preparation of an injection containing 500 mg lepidopteran, 250 mg glucose and injection water to 5 ml.\nSeveral analogs of naturally-occurring cecropins, sarcotoxins and lepidopterans have been reported. For example, it is reported in Andreu et al, Proc. Natl. Acad. Sci. USA, vol. 80, pp. 6475-6479 (1983) that changes in either end of the amino acid sequence of cecropin generally result in losses in bactericidal activity in varying degrees against different bacteria. It is reported in Andreu et al (1985) mentioned above that Trp.sup.2 is clearly important for bactericidal activity of cecropin, and that other changes in the 4, 6 or 8 position have different effects on different bacteria. From the data given in Table II at page 1687 of Andreu et al (1985), it appears that almost any change from natural cecropin generally adversely affects its bactericidal activity. Cecropin is defined in International Publication WO86/04356 to include bactericidally active polypeptides from any insect species and analogs, homologs, mutants, isomers and derivatives thereof having bactericidal activity from 1% of the naturally-occurring polypeptides up to 100 times or higher activity of the naturally-occurring cecropin. Other references generally discuss the effects of the .alpha.-helix conformation and the amphiphilic nature of cecropin and other lytic peptides.\nIt is known that lysozyme and attacins also occur in insect hemolymph. For example, it is reported in Okada et al, Biochem. J., vol. 229, pp. 453-458 (1985) that lysozyme participates with sarcotoxin against bacteria, but that the bactericidal actions are diverse. Steiner et al mentioned above suggests that lysozyme plays no role in the antibacterial activity of insect hemolymph other than to remove debris following lysis of bacteria by cecropin. Merrifield et al, Biochemistry, vol. 21, pp. 5020-5031 (1982) and Andreu et al (1983) mentioned above state that cecropin purified from insect hemolymph may be contaminated with lysozyme, but demonstrate that the synthetically prepared cecropin is as bactericidally active as purified cecropin from insect hemolymph.\nThe treatment of eucaryotic pathogens and other cells with lytic peptides, and novel lytic peptides, is the subject matter of U.S. Ser. No. 102,175, filed Sep. 29, 1987, now abandoned by Jaynes, Enright, White and Jeffers, which is hereby incorporated herein by reference.\nApproximately 70% of the world's human population lives in underdeveloped countries, and have diets nutritionally inadequate in proteins, fats and calories. Malnutrition in these countries is wide-spread and persistent. Protein malnutrition can usually be attributed to a deficiency in the diet of one or more of the essential amino acids. The major food staples of many underdeveloped countries, cereals and tubers, are deficient in most limiting essential amino acids. When a major portion of the diet consists of such staples, the result is limiting essential amino acid deficiency. In children, this condition is particularly debilitating, because of the large requirement for high quality polypeptide needed for normal growth and development.\nProtein malnutrition of this type could be alleviated or eliminated by adding supplements to the diet, which supplements contained polypeptides high in these limiting essential amino acids. Such polypeptides would have to be susceptible to digestion by normal human or animal proteases. Further, such polypeptides would be manufactured by cloning and expression of synthetic DNA.\nSynthetic DNA of a desired sequence can now be constructed using modern chemical techniques, and the DNA can then be cloned into various microorganisms using recombinant DNA technology. It is known, for example, that a synthetic DNA which codes for poly(1-aspartyl-1-phenylalanine) can be cloned and expressed in E. coli. Doel et al, Nucleic Acids Research, Vol. 8, No. 20, pp. 4575-4592 (1980). Also Tangus et al, Applied and Environmental Microbiology, Vol. 43, No. 3, pp. 629-635 (March, 1982), have obtained expression of a cloned homopolymeric synthetic DNA sequence coding for poly-1-proline.\nA number of United States patents have issued concerning the synthesis of polypeptides by conventional polypeptide sequencing, or by using DNA technology. Such patents include U.S. Pat. No. 3,796,631 to Choay et al, U.S. Pat. No. 3,850,749 to Kaufman et al, U.S. Pat. No. 3,299,043 to Schramm et al, U.S. Pat. No. 3,594,278 to Naylor, and U.S. Pat. No. 3,300,469 to Bernardi et al. U.S. Pat. No. 4,338,397, issued Jul. 1982 to Gilbert et al, discloses a method for synthesizing within a bacterial host, and secreting through the membrane of the host, a selected mature protein or polypeptide using these DNA techniques. The polypeptide can be any selected polypeptide such as proinsulin, serum albumin, and the like.\nPeptides high in limiting essential amino acids and transformed plants expressing the same are the subject matter of U.S. Ser. No. 837,722, filed Mar. 7, 1986, now abandoned by Jaynes; and U.S. Ser. No. 837,211 filed Mar. 8, 1986, now abandoned by Jaynes, which are hereby incorporated herein by reference."}
{"text": "This invention relates generally to turbine engine stator assemblies, and more particularly, to apparatus and method for controlling operating clearance between a stationary shroud surface in a turbine engine stator assembly and a rotating surface of juxtaposed blading members.\nForms of an axial flow turbine engine, typically a gas turbine engine, include rotating assemblies radially within stationary assemblies that assist in defining a flowpath of the engine. Examples include a rotary compressor assembly that compresses incoming air, and a rotary turbine assembly that extracts power from products of engine fuel combustion. Such assemblies comprise stages of rotating blades within a surrounding stator assembly that includes a shroud surface spaced apart from cooperating surfaces of the rotating blades. Efficiency of a turbine engine depends, at least in part, on the clearance or gap between the juxtaposed shroud surface and the rotating blades. If the clearance is excessive, undesirable leakage of engine flowpath fluid will occur between such gap resulting in reduced engine efficiency. If the clearance is too small, interference can occur between the rotating and stationary members of such assemblies, resulting in damage to one or more of such cooperating surfaces.\nComplicating clearance problems in such apparatus is the well known fact that clearance between such turbine engine assemblies changes with engine operating conditions such as acceleration, deceleration, or other changing thermal or centrifugal force conditions experienced by the cooperating members during engine operation. Clearance control mechanisms for such assemblies, sometimes referred to as active clearance control systems, have included mechanical systems or systems based on thermal expansion and contraction characteristics of materials for the purpose of maintaining selected clearance conditions during engine operation. Such systems generally require use of substantial amounts of air for heating or cooling at the expense of such air otherwise being used in the engine operating cycle. Provision of an improved means for active clearance control that reduces the need for engine flowpath fluid for such heating or cooling could enhance engine efficiency."}
{"text": "1. Field\nThe embodiments described below relate generally to the delivery of radiation therapy to a patient according to a “dynamic strobe” delivery scheme. In some embodiments, the “dynamic strobe” delivery scheme may encompass and/or seamlessly combine with one or more other radiation therapy delivery methods.\n2. Description\nAccording to conventional radiation therapy, a beam of radiation is directed toward a tumor located within a patient. The radiation beam delivers a predetermined dose of therapeutic radiation to the tumor according to a treatment plan. The delivered radiation kills cells of the tumor by causing ionizations within the cells. A major concern is limiting the damage to healthy tissue surrounding the tumor.\nFIG. 1 illustrates a conventional patient treatment process that includes radiation therapy. According to some examples of process 100, image data of a patient is acquired, and a target volume and critical internal structures are identified based on the image data during diagnosis (105). A radiation dose is prescribed (110) for achieving desired results with respect to the target volume while minimizing damage to the critical structures. Next, a treatment plan for delivering the dose is determined (115).\nThe treatment plan is then delivered (125) to the patient during several sessions, or “fractions”, spaced over some period of days. Prior to each fraction, the patient is positioned (120) as required by the treatment plan. Such positioning may involve the use of lasers, skin markers, etc.\nVarious methods or modes of radiation therapy delivery have been proposed and utilized. A radiation therapy delivery system typically includes hardware and a control system optimized for one specific mode of radiation therapy delivery. Some modes of radiation therapy delivery include, for example, conventional IMRT, Dynamic Modulated Arc Therapy, CT-Guided IMRT (Intensity Modulated Radiation Therapy), and Volumetric Modulated Arc Therapy.\nIn a conventional IMRT system, beams with modulated intensity are generated at a number of fixed positions or angles around the patient. These beams are then delivered with the gantry stationary at each fixed position. The beam intensity is modulated by either superimposing several shapes at a fixed position (“Step and Shoot IMRT”) or by moving leaves of a multileaf collimator (“MLC”) across the beam with varying speeds (“Sliding Window IMRT”). In a Dynamic Modulated Arc therapy system the gantry of the delivery system performs a contiguous rotational motion (360 degrees or less per arc). Throughout the rotational motion, the beam remains on at constant dose rate, and the MLC leaves constantly re-form to maintain a shape which is conformal with the shape of the tumor, as viewed from the respective angle. The throughput of Dynamic Modulated Arc therapy is greater than in conventional IMRT, but requires more sophisticated control of the MLC leaves.\nIn a CT-Guided IMRT system, such as a TomoTherapy® system provided TomoTherapy Incorporated, a linear accelerator is mounted in a ring-shaped gantry and moves in a 360 degree rotation around the patient. The beam is always on during this motion and is partially blocked or unblocked by rapidly opening and closing the MLC leaves as the gantry rotates. In a Volumetric Modulated Arc Therapy system, such as provided by the RapidArc™ radiotherapy technology from Varian Medical Systems, an L-shaped gantry performs a 360 degree rotation around the patient. The beam is constantly on and the dose rate may be modulated. Also, the MLC leaves are in constant motion, thus creating different shapes as the gantry rotates. These latter two systems require complex and expensive components for synchronizing and controlling motion of the gantry and the MLC leaves.\nDuring the treatment planning stage (115), a decision is made regarding which delivery mode to use. The decision may be based on a number of factors, including for example the patient's diagnosis, delivery system constraints, time issues related to schedules and availability of the patient and/or radiation therapy systems, etc. After the specific mode of delivery is decided, a treatment plan that is adapted to the delivery mode is determined.\nSelection of a delivery mode involves trade-offs or compromises as described above. Other radiation therapy delivery modes which may be particularly suitable to certain applications are desired."}
{"text": "Hitherto, in order to manufacture a surface magnet rotary electric machine at low cost, a combination of a rotor core having a perfectly circular cross section and a plate-shaped permanent magnet having a flat bonding surface has been used. It is difficult to bond the permanent magnet having the flat bonding surface to an outer peripheral surface of the rotor core having the perfectly circular cross section with highly accurate positioning of the permanent magnet. Therefore, retaining portions each having a polygonal columnar shape are provided to ends of the rotor core in an axial direction of the rotor core (see, for example, Patent Literature 1)."}
{"text": "1. Field of the Invention\nThe present invention relates to photovoltaic element useful as a solar cell or the like, the structure of an electrode thereof, and a process for producing the photovoltaic device.\n2. Related Background Art\nIn recent years, the global temperature rise caused by increase of carbon dioxide in the atmosphere, namely the greenhouse effect, has become a great problem. Therefore, clean energy sources are increasingly demanded. Nuclear power generation which does not cause the greenhouse effect involves the problem of disposal of radioactive waste. A safer and cleaner energy source is thus required.\nAmong promising clean energy sources, solar cells are attracting particular attention because of the safety, cleanness, and ease of handling thereof.\nOf the solar cells, amorphous silicon type solar cells which employ an amorphous silicon semiconductor are promising because they have advantages of ease of large area cell production, high light absorbance enabling thin film operation, and so forth, although the photovoltaic conversion efficiency is lower than that of crystalline silicon type solar cells.\nFIGS. 24A to 24D illustrate schematically a conventional photovoltaic element for comparison with the one of the present invention. FIGS. 24A to 24C illustrate the steps of forming a collecting electrode, and FIG. 24D is a cross-sectional view along the line 24D--24D in FIG. 24C. The solar cell element 500 of FIGS. 24A to 24D is prepared by successively laminating a lower electrode layer 502 on the surface of a substrate 501, a semiconductor layer 503 thereon, and an upper electrode layer 504 further thereon.\nIn such a solar cell element, in order to completely electrically separate the upper electrode layer from the lower electrode layer, a part 505 of the upper electrode layer is removed, and collecting electrodes 506 for the upper electrode layer 504 are provided on the surface of the upper electrode layer (FIG. 24A). In one method, for example, the collecting electrodes 506 are prepared by applying an electroconductive paste on the face of the photovoltaic laminate by screen printing and heat curing it. By this method, electrodes having a line width of 100 to 150 .mu.m, and thickness of 10 to 20 .mu.m are obtainable industrially. The electroconductive paste includes various materials. For amorphous silicon type solar cells, for which high temperature treatment is not suitable, polymer pastes comprising a thermosetting resin such as of polyester, epoxy and phenolic type, and a fine particulate material such as silver and copper dispersed therein are frequently used.\nOn the collecting electrodes 506 a bus bar electrode 507 is provided which further collects the generated power from the above collecting electrode 506 (FIG. 24B). Then crossing points of the bus bar electrode 507 over the collecting electrodes 506 are connected by applying an electroconductive adhesive 508 in spots and curing it in an air drier (FIG. 24C). Thus a leadout electrode which outputs the power from the upper electrode 504 is prepared by electrically connecting the collecting electrodes 506 with the bus bar electrode 507. Insulation tapes 509 are applied at the ends of the solar cell element 500 to ensure electrical separation of the bus bar electrode 507 from the substrate.\nIn producing this type of element, three steps are necessary: (1) registration of bus bar electrode to a prescribed position, (2) application of an electroconductive adhesive in spots at prescribed positions with the bus bar electrode fixed, and (3) curing the electroconductive adhesive by heating in an air drier or an IR oven. The steps involve many working operations and take a long time, and thus are disadvantageously not suitable for mass production. In the example shown in FIGS. 24A and 24B, the problems are not so serious since the area of the photovoltaic element is small and the number of the collecting electrodes is small. However, in the production of a photovoltaic element of a larger area, the number of the collecting electrodes is larger and a plurality of the bus bar electrodes are necessarily employed. In such a case, the electrically connected points increase in number and the working time becomes correspondingly longer, which is undesirable in respect of productivity.\nUsual family consumption of electric power is about 3 KW per family. To supply 3 KW of power by means of solar cells, the solar cell needs to have a light-receiving area of as large as 30 m.sup.2, by assuming the photovoltaic conversion efficiency of the cell to be 10%. Such a large solar cell is required to have a bus bar to collect the generated power, which increases the cell element production steps. The larger the number of the production steps, and the larger the area of the cell, the more the defects of the element are developed. The defects cause shunting and short circuits which lower the photovoltaic conversion efficiency. If the defect is distant from the electrodes or the bus bar, loss of current is relatively small because of high resistance to the current flowing into the defective portion. On the other hand, if the defect is beneath the electrode or the bus bar, the loss of current is large.\nTo solve these problems, electrode constructions are disclosed which are suitable for a large-area solar cell without using the bus bar. For example, U.S. Pat. No. 4,260,429 discloses a process in which a copper wire is covered with a solid polymer containing electroconductive particles and is attached as the electrode to a solar cell. U.S. Pat. No. 5,084,107 discloses a process in which a metal wire is connected and fixed by an electroconductive adhesive to a surface of a photovoltaic element. In these methods, the electrode is formed by covering an electroconductive wire with an electroconductive particle-containing solid polymer (electroconductive adhesive) with a low ohmic loss even with an electrode length of 10 cm or more.\nHowever, in a study of electrode constitution and reliability of solar cells conducted by the inventors of the present invention, it was found that the electrodes formed by the methods of the above U.S. Patents are insufficient in adhesion at the interface between the electroconductive wire or the metal member and the electroconductive particle-containing solid polymer or the electroconductive adhesive, and the electrodes are not uniform in the width or diameter. The insufficient adhesion at the interface between the electroconductive wire or the metal member and the electroconductive adhesive causes initial power loss, increase of series resistance by peeling at the interface in a long-term run, resulting in a drop in conversion efficiency, and other problems in reliability. Further, in the above elements, shunting and low yield may be caused, depending on the resistance of the electrode layer formed from the electroconductive adhesive, and disadvantageous migration of ionic substances may be caused by interaction with water, resulting in leakage due to humidity in practical use.\nIn the above U.S. Patents, the coated wires are bonded by heating or pressing, but neither the apparatus nor the method are shown specifically. However, it has been found that the usual method of heat-press bonding causes spreading of the coating polymer, which increases shadow loss, and particularly in a large-area photovoltaic element, the pressure applied to the coated wire becomes non-uniform to give non-bonded portions which disadvantageously increase the series resistance.\nElectrical connection is now further considered for outputting the generated power of the photovoltaic element.\nFIG. 17 is a schematic plan view of a photovoltaic element from the front side (light-receiving side) for comparison. The photovoltaic element of FIG. 17 is constituted by an electroconductive substrate for supporting the entire photovoltaic element, a non-crystalline semiconductor layer, an electrode layer, collecting electrodes, and lead-out terminals successively formed over the substrate. The electroconductive substrate is made of a metallic material such as stainless steel. The semiconductor layer is constituted by a back-face reflection layer, a p-type semiconductor layer, an i-type semiconductor layer, and an n-type semiconductor layer successively arranged from the bottom layer. These semiconductor layers are formed and laminated by a film-forming method such as CVD (chemical vapor deposition) so that the light may be efficiently absorbed and converted to electric power. The aforementioned electrode layer is a light-transmissive electroconductive film made of indium oxide or the like, serving both as a reflection-preventing means and as a power-collecting means.\nThe light-transmissive electroconductive film is formed by application of an etching-paste containing FeCl.sub.3, and AlCl.sub.3, or the like by screen printing, and subsequent heating. An etched groove 17401 is formed by removing the light-transmissive electroconductive film by etching in a line shape. This partial removal of the light-transmissive electroconductive film is conducted for preventing short circuits between the substrate and the light-transmissive electroconductive film on the effective light-receiving area of the photovoltaic element. Such short-circuits may occur in cutting of the outer periphery of the photovoltaic element.\nOn the surface of the above photovoltaic element, collecting electrodes 17402 are formed for efficiently collecting the generated power. The collecting electrodes 17402 are formed by using a fine metal wire of low resistance such as copper as the core material, applying an electroconductive adhesive on the outer surface of the metal wire for adhesion, drying the adhesive, cutting the wire in a predetermined length, arranging the cut wires, and heat-bonding the wires on the surface of the effective light-receiving area by hot pressing.\nThe power collected by the collecting electrode 17402 is transmitted to lead-out terminals 17403 provided on both ends of the element. The lead-out terminals 17403 are foils made of a low-resistance metal such as copper, with an insulating member 17404 as the lowest layer to insulate the foils from the surface of the photovoltaic element.\nThe connection between the collecting electrodes 17402 and the lead-out terminals 17403 is conducted by spot-like application of an electroconductive adhesive 17405 to ensure reliable connection.\nThis process involves the problems of: necessity of spot-like application of an electroconductive adhesive, requiring a heat-curing step for curing the electroconductive adhesive, thereby increasing the working steps, taking a long time, and requiring a complicated apparatus for conducting these steps; the high cost of the electroconductive adhesive; the protrusion of the spot-wise applied electroconductive adhesive requiring a thick surface coating material, thus raising the production cost; deterioration of the surface of the terminal material such as copper by oxidation and other causes during the heating process of the element prior to the application of the electroconductive adhesive, whereby a sufficiently low connection resistance is not obtained by the application of the electroconductive adhesive."}
{"text": "1. Field of the Invention\nThis invention relates generally to a vehicle rollover avoidance system and, more particularly, to a vehicle rollover avoidance system that employs a roll control factor and a yaw rate stability control factor to control semi-active suspension dampers to mitigate the risk of vehicle rollover.\n2. Discussion of the Related Art\nIt is known in the art to mitigate a potential vehicle rollover using differential braking control, rear-wheel steering control, front-wheel steering control, or any combination thereof. A vehicle rollover avoidance system may receive vehicle dynamics information from various sensors, such as yaw rate sensors, lateral acceleration sensors and roll rate sensors, to determine the proper amount of action to be taken to detect a potential vehicle rollover. A balance typically needs to be provided between estimating the vehicle roll motion and the vehicle yaw motion to provide the optimal vehicle response. Thus, it is usually necessary to detect certain vehicle conditions to provide the roll detection. To precisely identify vehicle roll stability conditions, it may be advantageous to know the vehicle's roll rate and roll angle because they are the most important states in vehicle roll dynamics.\nUnder normal driving conditions, drivers can direct the vehicle to the desired heading through the control of the steering wheel. When the vehicle is turning, there are actually three motions taking place with the vehicle. Particularly, a turning motion, or yaw, is occurring, as the vehicle body is turning around an imaginary access vertical to the ground through the so-called yaw-center of the vehicle. Also, there is subtle vehicle sliding laterally, sometimes in the direction of the turn and sometimes away from the turn, depending mainly on the vehicle speed. Further, a tilting motion or roll motion occurs as the vehicle's body is turning around an imaginary axis parallel to the ground through the so-called roll-axis of the vehicle.\nUnder normal vehicle maneuvering conditions, the tire/road contact surfaces can generate sufficient forces to sustain the desired vehicle motions, and drivers are accustomed with these motions as they occur. However, when the vehicle maneuver starts approaching limit-handling conditions, the tire/road contact surfaces can no longer sustain the desired yaw motion and side-slip motion, and the vehicle body will exhibit an increased roll motion. As a result, a discrepancy will build up between the vehicle's yaw rate and its desired yaw rate, and between the vehicle's side-slip velocity and its desired side-slip velocity. Further, if the roll motion becomes too large, the vehicle may roll over."}
{"text": "This invention relates to high energy electrochemical power cells. More particularly, it relates to cells having an oxidizable anode material and a liquid active cathode material and a solid current collector.\nRecently there has been a rapid growth in portable electronic products requiring electrochemical cells to supply the energy. Examples are calculators, cameras and digital watches. These products, however, have highlighted the deficiencies of the existing power cells for demanding applications. For example, digital watches were developed using the silver oxide cell, and although these watches have become popular, it is now generally recognized that the least developed component of the digital watch system is the power cell. In particular, the energy density of the silver cells is such that thin, stylish watches with reasonable operating life are difficult to make. Additionally, these cells have poor storage characteristics, low cell voltages, and leakage problems.\nIn an effort to develop a cell that addresses one or more of the foregoing problems, substantial work has been done with cell chemistries using a lithium anode. The cathode and electrolyte material consisting of a solvent and solute vary. Indeed the literature is replete with examples of lithium anode cells with different cathodes and electrolytes. The electrical characteristics of these cells such as energy per unit volume, called energy density; cell voltage; and internal impedance vary greatly.\nAmong all of the known combinations of lithium anodes with different cathodes and electrolytes, those believed to have among the highest energy density and lowest internal impedance use certain inorganic liquids as the active cathode depolarizer. This type of cell chemistry is commonly referred to as \"liquid cathode\".\nEarly liquid cathode cells used liquid sulfur dioxide as the active cathode depolarizer as described in U.S. Pat. No. 3,567,515 issued to Maricle, et. al. on Mar. 2, 1971. Since sulfur dioxide is not a liquid at room temperature and at atmospheric pressure, it proved to be quite a difficult chemistry with which to work. More importantly, sulfur dioxide cells are unsafe for most consumer applications due to their propensity to explode under certain circumstances.\nA major step forward in the development of liquid cathode cells was the discovery of a class of inorganic materials, generally called oxyhalides, that are liquids at room temperature and also perform the function of being the active cathode depolarizer. Additionally, these materials may also be used as the electrolyte solvent. Liquid cathode cells using oxyhalides are described in U.S. Pat. No. 3,926,669 issued to Auborn on Dec. 16, 1975, and in British Pat. No. 1,409,307 issued to Blomgren et al. on Oct. 18, 1975. At least one of the oxyhalides, thionyl chloride (SOCl.sub.2), in addition to having the general characteristics described above, also provides substantial additional energy density.\nAs described in the Auborn and Blomgren patents, the anode is lithium metal or alloys of lithium and the electrolyte solution is an ionically conductive solute dissolved in a solvent that is also an active cathode depolarizer.\nThe solute may be a simple or double salt which will produce an ionically conductive solution when dissolved in the solvent. Preferred solutes are complexes of inorganic or organic Lewis acids and inorganic ionizable salts. The requirements for utility are that the salt, whether simple or complex, be compatible with the solvent being employed and that it yield a solution which is ionically conductive. According to the Lewis or electronic concept of acids and bases, many substances which contain no active hydrogen can act as acids or acceptors or electron doublets. In U.S. Pat. No. 3,542,602 it is suggested that the complex or double salt formed between a Lewis acid and an ionizable salt yields an entity which is more stable that either of the components alone.\nTypical Lewis acids suitable for use in the present invention include aluminum chloride, antimony pentachloride, zirconium tetrachloride, phosphorus pentachloride, boron fluoride, boron chloride and boron bromide.\nIonizable salts useful in combination with the Lewis acids include lithium fluoride, lithium chloride, lithium bromide, lithium sulfide, sodium fluoride, sodium chloride, sodium bromide, potassium fluoride, potassium chloride and potassium bromide.\nThe double salts formed by a Lewis acid and an inorganic ionizable salt may be used as such or the individual components may be added to the solvent separately to form the salt. One such double salt, for example, is that formed by the combination of aluminum chloride and lithium chloride to yield lithium aluminum tetrachloride.\nIn addition to an anode, active cathode depolarizer and ionically conductive electrolyte, these cells require a current collector.\nAccording to Blomgren, any compatible solid, which is substantially electrically conductive and inert in the cell, will be useful as a cathode collector since the function of the collector is to permit external electrical contact to be made with the active cathode material. It is desirable to have as much surface contact as possible between the liquid cathode and the current collector. Therefore, a porous material is preferred since it will provide a high surface area interface with the liquid cathode material. The current collector may be metallic and may be present in any physical form such as metallic film, screen or a pressed powder. Examples of some suitable metal current collectors are provided in Table II of the Auborn Patent. The current collector may also be made partly or completely of carbon. Examples provided in the Blomgren Patent use graphite.\nElectrical separation of current collector and anode is required to insure that cathode or anode reactions do not occur unless electrical current flows through an external circuit. Since the current collector is insoluble in the electrolyte and the anode does not react spontaneously with the electrolyte, a mechanical separator may be used. Materials useful for this function are described in the Auborn Patent.\nAlthough the varied cells described in the Blomgren and Auborn patents may be feasible, much of the recent interest is in cells using thionyl chloride as the active cathode depolarizer and electrolyte solvent. This results from thionyl chloride's apparent ability to provide greater energy density and current delivery capability than other oxyhalide systems. Yet even though thionyl chloride cells have proven to be the best performer among the oxyhalides, they have not lived up to expectations on energy density or internal impedance. Furthermore, the thionyl chloride cell is equally if not more dangerous than the sulfur dioxide cell. As a result, all known efforts to commercialize cells using this chemistry have failed."}
{"text": "1. Field of the Invention\nThe present invention relates to an imaging device which uses a solid state image pickup element. More particularly, the present invention relates to an imaging device which (a) adjusts the ratio of color excitation values (for example, red/green/blue \"RGB\") of the image pickup element in correspondence with the exit pupil position of an associated optical unit, and (b) adjusts the ratio of the color excitation values in correspondence with the size of an aperture controlling light passing through the associated optical unit.\n2. Description of the Related Art\nA conventional imaging device converts optical images to image signals using a solid state image element, such as a CCD element (charge coupled device) or a MOS element (metal oxide semiconductor element). Moreover, recent technological developments have attempted to produce such an imaging device with a more compact solid state image element and to increase the number of pixels in the imaging device. Unfortunately, as an imaging device is made with a more compact solid state image element, the aperture efficiency of the image pickup element is reduced and the signal-to-noise (S/N) ratio of the image signals is decreased.\nFIG. 1 is a diagram illustrating a cross section of a conventional CCD image pickup element. Referring now to FIG. 1, an image pickup element 1 includes a light receiving unit 1a (such as a \"pixel\") which converts light to an electric charge, and a transfer unit 1b which transfers the electric charge from light receiving unit 1a. Light receiving part 1a and transfer unit 1b are formed on the surface of image pickup element 1. An on-chip micro-lens 2 forms a condenser lens which corresponds to light receiving units 1a. On-chip micro-lens 2 is arranged on the surface of image pickup element 1.\nWith image pickup element 1, the light incident on-chip micro-lens 2 is focused on light receiving unit 1a. For example, as illustrated in FIG. 1, a light ray A and a light ray B are both focused by on-chip micro-lens 2 so that the respective light rays strike light receiving unit 1a. Therefore, the amount of light received by light receiving unit 1a is increased. Consequently, the level of the signals that are optoelectrically converted by light receiving unit 1a becomes larger, and the imaging device can output image signals with a high S/N ratio.\nFIG. 2 is a diagram illustrating a conventional electronic still camera using a conventional CCD image pickup element. Referring now to FIG. 2, an aperture 4 and a mirror 5 are arranged in the optical axis of a photographic lens 3, and image pickup element 1 is arranged on the focal plane of photographic lens 3. On-chip micro-lens 2 is formed on the light receiving surface of image pickup element 1. Image pickup element 1 produces three excitation values (referred to as an \"R output\", a \"G output\" and a \"B\" output). A signal processing unit 7a processes the R output, G output and B output of image pickup element 1, and produces corresponding image signals. The G output of image pickup element 1 is connected \"as is\" (that is, \"directly\") to signal processing unit 7a. The R output and the B output of image pickup element 1 are connected to signal processing unit 7a through variable gain amplifiers 6a and 6b, respectively. A recording unit 7 records the image signals produced by signal processing unit 7a.\nA light measurement unit 8a measures the subject brightness (that is, the brightness of a subject (not illustrated)) and is arranged in a position which is illuminated by the light reflected from mirror 5. An exposure calculation unit 8 is connected to light measurement unit 8a. Also, control terminals (not illustrated) of aperture 4 and image pickup element 1 are individually connected to output terminals of exposure calculation unit 8. A color measurement unit 9a measures the color of the ambient light. A white balance control unit 9 is connected to the control terminals of variable gain amplifiers 6a and 6b, and to color measurement unit 9a. In addition, a control unit 10 controls, and is connected to, photographic lens 3, image pickup element 1, signal processing unit 7a, recording unit 7, exposure calculation unit 8, and white balance control unit 9. A release button 10a is also connected to control unit 10. Release button has a half-push position and a full-push position and is pushed by a photographer to either the half-push position or the full-push position to initiate specific camera operations.\nIn an electronic still camera as illustrated in FIG. 2, when release button 10a is pushed to the half-push position, exposure calculation unit 8 incorporates the light measurement value of the subject brightness based on a measurement by light measurement unit 8a, and calculates the correct aperture value and exposure. Moreover, white balance control unit 9 incorporates the color measurement value of the ambient light based on a measurement by color measurement unit 9a, and controls the gain of variable gain amplifiers 6a and 6b in correspondence with the ratio of the three excitation values (RGB) of the ambient light. In this state, if release button 10a is pushed to the full-push position, mirror 5 flies up, and exposure calculation unit 8 adjusts aperture 4 to the correct aperture value.\nThe amount of light incident on image pickup element 1 from photographic lens 3 is restricted by aperture 4, and an optical image is focused on the light receiving plane of image pickup element 1. On-chip micro-lens 2 allows image pickup element 1 to raise the light receiving efficiency, and to produce image signals with a high S/N ratio. The white balance of the image signal produced by image pickup lens 1 is adjusted by variable gain amplifiers 6a and 6b. Signal processing unit 7a processes the image signal produced by image pickup lens 1 and adjusted by variable gain amplifiers 6a and 6b. Such processing performed by signal processing unit 7a can include, for example, gamma correction, and gain adjustment. The processed image signal produced by signal processing unit 7a is recorded in record unit 7.\nAs the exit pupil position of photographic lens 3 approaches image pickup element 1, the optical rays incident on on-chip micro-lens 2 from the side direction become more numerous. Because the light rays incident from the side direction are strongly affected by color aberrations on the axis of on-chip micro-lens 2, the size of spots focused on light receiving unit 1a change for every wavelength of light. For this reason, the amount of light output from light receiving unit 1a varies for every wavelength of light. This variation in the amount of light output from light receiving unit 1a causes the color phase of the light output from light receiving unit 1a to vary.\nFIG. 3(A) is a graph indicating the output of an image pickup element versus the distance between the exit pupil position of an associated optical system and the image pickup element. As illustrated by FIG. 3(A), the ratio of the R output of an image pickup element to the G output of the image pickup element becomes large as the exit pupil position approaches (becomes \"near\") the image pickup element. As a result, an image produced by the image pickup element appears reddish as the exit pupil position approaches the image pickup element. Moreover, the ratio of B output of the image pickup element to the G output of the image pickup element becomes smaller as the exit pupil position approaches the image pickup element. As a result, the blueness of an image produced by the image pickup element becomes thin as the exit pupil position approaches the image pickup element.\nIn particular, if a zoom lens is used with an image pickup element, the exit pupil position is moved greatly forward and backward following the adjustment of the image angle. Therefore, there is a wide range of fluctuations in the ratio of RGB output from the image pickup element, and the color phase of the image varies greatly. This kind of color phase variation cannot be corrected by an external light white balance adjustment. Therefore, in a conventional imaging device, the color phase of an image produced by an image pickup element undesireably varies in correspondence with the exit pupil position of an associated photographic lens.\nMoreover, referring to the imaging device illustrated in FIG. 2, if aperture 4 is set to the open aperture side (that is, near the fully open position), the light rays incident from the side in relation to on-chip micro-lens 2 become numerous. Therefore, there are changes in the color phase of the image produced by the image pickup element.\nFIG. 3(B) is a graph indicating the output of an image pickup element versus the aperture value of an aperture controlling the light incident on the image pickup element. As illustrated by FIG. 3(B), the ratio of R output of an image pickup element to the G output of the image pickup element becomes large as the aperture approaches the open aperture side. As a result, the image produced by the image pickup element appears reddish as the aperture approaches the open aperture side. Moreover, the ratio of B output of the image pickup element to the G output of the image pickup element becomes small as the aperture approaches the open aperture side. As a result, the blueness of the image produced by the image pickup element becomes thin as the aperture approaches the open aperture side,\nMoreover, if the white balance is manually selected by the photographer, the white balance is often set to the lowest permissible limit. Therefore, the ratio of RGB output in the state of the open aperture changes, and there is an increased probability that the white balance will become unnatural.\nIn view of the above, it is difficult to accurately reproduce the color phases of the subject because the ratio of RGB output of an image pickup element changes depending on the aperture value."}
{"text": "1. Field of the Invention\nThis invention relates generally to tools for manipulating pipe, and in particular, to hand-held tools for engaging various portions of irrigation pipe to manipulate the pipe without undue physical strain on the operator.\n2. Description of the Prior Art\nThere are many types of irrigation systems known in the prior art. One common system with which the present invention is particularly well suited for use includes a series of pipe joints connected together across an end of a field. The pipe joints each typically have a female bell connector at one end and a male connector at the other end. Several (e.g., 6 to 18) adjustable gate valves are typically spaced along the length of each pipe joint. The pipe joints are typically made of aluminum or plastic and range in size from 20 to 36 feet in length and 6 to 12 inches in diameter, with a 30 foot length and 8 to 10 inch diameter being the most common. The pipe joints typically weigh approximately 50 to 150 pounds, depending on the size and material of the pipe and the amount of sediment accumulated within the pipe.\nIrrigating is often the most labor intensive task on an agricultural farm in areas having inadequate rainfall. The irrigation system described above using pipe joints requires the pipe joints to be carried to the field to be irrigated at the beginning of each irrigating season, laid across the end of the field, connected together by inserting the male ends into the female ends of adjacent pipes, and adjusting the flow of water through the adjustable gate valves each time water is applied to the crop. It is often necessary to rotate the pipe after it has been connected together across the end of the field in order to direct the flow of water out of the gate valves to an appropriate angle away from the pipe. In addition, at the end of the season, the pipe usually must be disconnected, picked up, and carried out of the field.\nTwo of the most physically demanding tasks with the irrigation pipe system described above are turning the pipe to adjust the angle of the water flow after the pipe joints are connected together (especially if the pipe is filled with water), and disconnecting the pipe joints at the end of the irrigating season. The difficulty of these tasks is increased after the pipe settles during the irrigation season (typically one to three months), after sediment accumulates within the pipe, or as vegetative growth is allowed to grow under and around the pipe. A rubber gasket is typically used to seal the connection joints against leakage, thus increasing the force required to rotate and disconnect the pipe joints.\nOne known prior art system for rotating irrigation pipe has a pair of clamping jaws and a scissor-type handle for engaging the pipe. This system requires the use of two hands to operate the handle to apply a clamping force to the pipe. Thus, in situations where the operator's hands are full of other tools, such as shovels and gate changer tools, the known turning tool is unwieldy and not very useful. This known tool is also relatively complex, making it more difficult and expensive to manufacture.\nThe common method of disconnecting the pipe joints at the end of the irrigation season is to grip the bell end of the pipe with the operator's hands (preferably wearing gloves to avoid cuts and abrasions) and twisting and pulling until the adjacent pipe connection comes apart. This method is, of course, very strenuous on the irrigator, often resulting in lower back pain, sore hands, and physical exhaustion. In addition, the irrigator is limited as to the amount of force he can apply without his hands slipping from the pipe.\nThus, there is a need for devices to reduce the physical exertion required for these tasks without the disadvantages of the prior art described above."}
{"text": "In client-server environments, remote desktop clients were designed to be used on client devices with display modules that are at least as equally large as display modules at the host server. Even if there is no display device at the host server, there is an assumption that the image that represents the data for display, generated at the client device, is generated for a display module that is of desktop size (e.g. >12″). However, challenges arise when the remote desktop client is installed at a client device having a display module that is smaller than the assumed size. For example, mobile devices generally have very small display modules (e.g. less than 3″), and generating an image to represent data meant for display on a much larger display module is challenging.\nOne solution is to enable a mobile device with a small display module to generate a portion of the data for display, in essence creating a peep bole into the larger host desktop window. Moving the peep hole around can be painful since navigating a large area with a trackball or thumbwheel is awkward. For example, a user may be reading a document (or contents of a window). The user starts at the top left corner and moves the peep hole from left to right for example by moving a cursor. At the end of the row/window the user must scroll across to the left margin and down one row/line. Getting to the new location is hence awkward.\nOne approach to this is provided in “Advancing interaction: ZoneZoom: Map navigation for smartphones with recursive view segmentation” by Robbins, D. C., Cutrell, E., Sarin, R., & Horvitz, E. (2004), and published in the proceedings of the working conference on advanced visual interfaces (AVI '04), (Gallipoli, May 2004), ACM press, 231-234. In this approach an information space is segmented into nine sub-segments, each of which is mapped to a key on the number keypad of a smartphone having a display module. The sub-segments can be chosen by the author of the information space or dynamically generated at run-time. To view a sub-segment, a user presses the appropriate button on the keypad to take advantage of “spring-loaded” view shifting, which allows users to jump between views of defined sub-segments. However, this approach is awkward if the user is unclear about what he/she wishes to view and may have to hunt between sub-segments to find the appropriate view."}
{"text": "The present invention relates to a method of manufacturing silver halide emulsion, particularly to a method of manufacturing a silver halide emulsion low in fog, excellent in graininess, low in aging fog, and excellent in aging stability.\nIn the area of silver halide photographic light-sensitive materials, various approaches are being made in pursuit of a much higher speed. As one of such approaches, there is proposed to use core/shell type silver halide grains different in chemical components or physical properties between the inner portion and the other portion, generally, silver halide grains different in silver halide composition (hereinafter referred to as a core/shell type emulsion). In using such core/shell type emulsions for the purpose of enhancing the sensitivity, it is proposed to make up the light-sensitive layer into a multilayered structure in which the upper layer (a layer on the light-irradiated side) is a high speed layer.\nIn general, core/shell type emulsion grains have a structure in which the inner silver halide composition of the grains differs from the silver halide composition of the other portion of the grains; and, present inventors have found the fact that a core/shell type emulsion comprising grains of this structure having a high iodide content phase (a portion in which the iodide content is higher than that in the other portion) gives a fog high than that comprising grains having no high iodide content phase. Such an increased fog causes a large deterioration in graininess and a substantial desensitization with respect to a sensitivity-fog ratio.\nParticularly, use of a core/shell type emulsion in a multilayered light-sensitive material involves risks to increase fog or aging fog and impair the shelf-life."}
{"text": "1. Field of the Invention\nThe present invention relates, generally, to medical devices. More particularly, the invention relates to implantable cardioverter defibrillators (ICDs).\n2. Background Information\nIn the past, various ICD devices and methods have been used or proposed. However, these devices and methods have significant limitations and shortcomings. Existing ICD's are primarily designed for chronic applications in that they produce enough shocks to treat a chronic condition wherein the patient is expected to have numerous episodes of sudden cardiac death (SCD) over an extended period. Typical defibrillators are capable of delivering from 150 to 350 full output shocks of 27 to 40 joules each (depending on the model and manufacturer). In order to deliver that many high energy shocks, the device must have sufficient battery capacity to cover the required delivered energy as well as system losses (about 30% is lost in the DC to DC converter). This necessarily adds bulk and weight to the device.\nPresently, most patients undergoing ICD implantation have exhibited at least one episode of fibrillation (SCI) and survived due for example to early CPR, trans-chest defibrillation and other care. Since one episode is typically a clear indication of high risk of having another one, an ICD is indicated. Other patients exhibit very early indications for being at high risk for SCD and an ICD is implanted prophylactically. Overall, about 40% of patients who have ICD's implanted do not have another episode during the next four years. However, these patients still need protection and typically another ICD must be implanted after the battery dies in three to five years, even though no shocks are delivered by the device.\nPatients who are not shocked by their ICD have unnecessarily had a large device capable of hundreds of shocks implanted. Large devices are uncomfortable and present an increased risk of infections, erosions, and certain psychological problems. A smaller device with a smaller battery (and possibly fewer functions) would serve these patients better. Such a device would have only a sufficient number of shocks available to save the patient from initial SCD episodes, whereupon the patient would immediately have a larger device with more shock capacity implanted. This type of device would be implanted in those patients who were considered at high risk, but have not yet had an episode (and may never have one) and in patients who have had a conventional large device which needs replacement, but who have not had a shock during the last several years. The concept of a prophylactic ICD is disclosed in U.S. Pat. No. 5,439,482.\nIt may seem obvious to a casual observer that to make a device with fewer shocks, one only need to use a smaller battery. That has not been the case, however, owing to other requirements of the battery. The battery must be capable of charging the output capacitor to its maximum output (27 to 40 joules) in a period of 6 to 10 seconds after detection of fibrillation. This typically requires from 0.7 to 1.0 amp of current during the charging period. With conventional batteries used in ICD's, Lithium Vanadium Pentoxide9 for example, the minimum size battery that meets the charging criteria has sufficient capacity for about 150 or more shocks. Thus, it has not been possible to make a limited shock device with a small battery. Virtually every battery's chemistry has this capacity/power relationship.\nAccordingly, it is an object of the present invention to provide an improved ICD which overcomes the limitations and shortcomings of the prior art, particularly those related to the limitations of prior art battery systems."}
{"text": "It is known to use power combiners comprising electrical conductors to bring together multiple high-frequency signal sources and/or to split a high-frequency signal. A power combiner that includes a second electrical conductor that is spaced apart from a first electrical conductor is known from EP 1 699 107 A1, in which the first electrical conductor is capacitively and inductively coupled to the second electrical conductor. Both the first electrical conductor and the second electrical conductor can include multiple windings to increase the inductive coupling between the conductors, in accordance with EP 1 699 107 A1.\nOther power combiners are known, for example, from U.S. Pat. No. 8,044,749 B1, DE 103 42 611 A1, and US 2014/0085019 A1."}
{"text": "Compound archery bows typically have a bowstring, on which an arrow may be nocked, along with one or more portions of cable other than the bowstring extending between the limbs of the bow. Such cable portions, sometimes referred to as “power cables”, are generally located at least partly within or close to an operating plane of the bowstring. The power cables thus interfere with shooting arrows.\nIn order to provide adequate room for the arrow, it is conventional practice to mount a cable guard on the bow to engage the central portions of the power cables and to displace them laterally a sufficient distance to one side of the operating plane of the bowstring to avoid interference with an arrow. One drawback associated with conventional cable guards is that, in displacing the center of a power cable laterally from its straight line position, they introduce a lateral component to the force exerted by the power cable against the limbs. This lateral torque not only decreases the accuracy of arrow flight, but also causes twisting of the limbs, cams, wheels and/or handle, and thereby contributes adversely to shortening their useful life. Conventional cable guards also cause the power cables to feed on and off of the cams and wheels at an angle. This may sometimes lead to the power cables becoming dislodged from the cams and/or wheels.\nThere exist a number of prior art systems, other than cable guards, for preventing the power cables from interfering with the shooting of arrows from compound bows. Examples include U.S. Pat. No. 5,623,915 to Kulacek and No. 6,729,320 to Terry, and U.S. patent application Ser. No. 11/968,459 to Evans.\nThe inventor has determined a need for further systems which do not require cable guards to prevent power cables from interfering with the flight of arrows."}
{"text": "Hydrocodone and hydromorphone are widely used semisynthetic narcotic analgesics possessing also useful antitussive properties. Hydrocodone is also an important intermediate for the synthesis of other opioid analgesics, e.g. dihydrocodeine.\nHydrocodone has been prepared by reduction of codeinone (Arch. Pharm (1920), 258, 295; J. Am Chem. Soc (1950), 72, 3247, U.S. Pat. No. 5,571,685) or thebaine (DE Pat. No. 441,613). However, since codeinone is obtained by oxidation of codeine in low to moderate yields, and thebaine is a relatively scarce alkaloid, these methods have little practical value.\nA number of methods describe convenient one-pot isomerization of codeine to hydrocodone and morphine to hydromorphone in the presence of noble metal catalysts (examples may be found in “The Chemistry of the Morphine Alkaloids” by Bentley, K. W.; U.S. Pat. No. 2,544,291; U.S. Pat. No. 5,847,142, WO Pat. No. 0134608). Unfortunately these processes require large amounts of expensive metals of the platinum group, or call for exotic organometallic catalysts. Also isomerization is complicated by uncontrollable side reactions, requiring introduction of laborious purification procedures, and resulting in substantial yield loss.\nHydrocodone and hydromorphone have been prepared by Oppenauer oxidation of dihydrocodeine or dihydromorphine correspondingly (U.S. Pat. No. 2,649,454; U.S. Pat. No. 2,654,756; Synth. Commun. (2000), 30, 3195). The starting materials for the oxidation, dihydrocodeine or dihydromorphine, may be synthesized by hydrogenation of morphine or codeine. Hydrogenation is typically performed in solvents, such as aqueous acetic acid (J. Org. Chem. (1950), 15, 1105), ethyl acetate (Synth. Commun. (2000), 30, 3195), or ethanol (Ann. (1923), 433, 269), that are not compatible with the reaction media for Oppenauer oxidation. This makes necessary solvent removal and isolation of the products of hydrogenation, and, compared to the convenient one-pot syntheses discussed above, adds additional steps to the process of hydrocodone or hydromorphone preparation.\nA need thus remains for an efficient, economical, and practical process for the production of hydrocodone, hydromorphone and related compounds."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the design of integrated circuits and more particularly to sense amplifiers.\n2. Description of the Background Art\nMany systems on an integrated circuit are designed to respond differently depending upon whether their input voltages are considered HIGH or LOW. Sometimes, an input voltage must be modified to conform to the HIGH or LOW state (e.g., during the period when the input voltage transitions between states). Sense amplifiers are circuits that detect a small voltage differential and increase or decrease the voltage to a level required by the system. An example of a system that utilizes sense amplifiers is a computer memory circuit. Information stored in the memory cells of a memory chip using sense amplifiers can be retrieved much faster than from a memory chip without sense amplifiers.\nAs shown in FIG. 1, a common static random access memory (SRAM) configuration generally designated 100 includes an array 105 of memory cells 110. Each memory cell 110 is connected to a word line 115, a bit line B 120, and a complement of the bit line, B 145. The memory cells 110 connected to each of the word lines 115 define a memory cell array row 125, and the memory cells connected to each of the bit line B 120 and a corresponding complement of the bit line B 145 define a memory cell array column 130. Each memory cell 110 stores information in the form of a voltage charge representing a logic value of LOW or HIGH. A voltage level equal to V.sub.DD represents the logic value of HIGH and V.sub.SS represents the logic value of LOW.\nBit lines B 120 and B 145 are connected to an equalization and precharge circuit 150. The precharge component of the equalization and precharge circuit 150 initially charges bit lines B 120 and B 145 to the voltage level of V.sub.DD. The equalization component of the equalization and precharge circuit 150 ensures that the voltages on bit lines B 120, .nu..sub.B, and B 145, .nu..sub.B, are initially at the same level.\nThe word lines 115 are connected to a row decoder 155. When a memory cell 110' is accessed, the row decoder 155 selects and changes the voltage of a word line 115' corresponding to memory cell 110'. A changed voltage signal (e.g., LOW to HIGH) from the word line 115' allows memory cell 110' to communicate with bits lines B 120' and B 145'. If memory cell 110' stores a logic value of HIGH, then .nu..sub.B will remain at HIGH and .nu..sub.B will decrease to LOW. If memory cell 110' stores a logic value of LOW, then .nu..sub.B will decrease to LOW and .nu..sub.B will remain at HIGH.\nBit lines B 120 and B 145 are connected to a sense amplifier 160 which detects and amplifies the difference in voltage between .nu..sub.B and .nu..sub.B. Depending on the difference between .nu..sub.B and .nu..sub.B, the sense amplifier 160 will output either V.sub.DD or V.sub.SS.\nConnected to the sense amplifier 160 is a column decoder 165. The column decoder 165, like the row decoder 155, includes a combination of logic circuits that select a logic signal from either one or a set of the memory cell array columns 130 for final output from SRAM 100.\nThe prior art described above suffers from a number of limitations. To store more information on a single memory chip, smaller memory cells are used. Smaller memory cells, however, use smaller transistors, which have less driving capability, resulting in a longer time for .nu..sub.B and .nu..sub.B to reach distinct HIGH or LOW voltage levels. To reduce the time required to read a memory cell, sense amplifiers are used to quickly detect the small voltage difference between .nu..sub.B and .nu..sub.B without having to wait for .nu..sub.B and .nu..sub.B to reach definite HIGH or LOW voltage levels. However, when .nu..sub.B and .nu..sub.B reach definite HIGH or LOW voltage levels before the operation of the sense amplifier, the operation of the sense amplifier is not required and consumes unnecessary power.\nWhat is needed is a sense amplifier design that overcomes the shortfalls of the sense amplifier designs known in the art."}
{"text": "The present invention relates to a programmable electronic calculator and, more particularly, to a programmable electronic calculator which includes a plurality of tape recorders as outer memories.\nA programmable electronic calculator or a portable, personal computer has been developed, which includes a tape recorder as an outer memory for storing a program or data. The programmable electronic calculator includes a control system for driving the tape recorder in order to conduct the read operation and/or the write operation of the program and the data into and from the tape recorder.\nHowever, in the conventional programmable electronic calculator, the programmable electronic calculator does not store the command for controlling the tape drive except for the read operation and/or the write operation of the program or the data.\nAccordingly, an object of the present invention is to provide a programmable electronic calculator which includes a control system for driving a tape recorder which is connected to the programmable electronic calculator as an outer memory.\nAnother object of the present invention is to provide a tape drive control system in a programmable electronic calculator, which independently drives a plurality of tape recorders connected to the programmable electronic calculator.\nOther objects and further scope of applicability of the present invention will become apparent from the detailed description given hereinafter. It should be understood however, that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description.\nTo achieve the above objects, pursuant to an embodiment of the present invention, at least two tape recorders are connected to the programmable electronic calculator.\nThe programmable electronic calculator stores commands for driving the tape recorders, which are applied to the tape recorders via an interface circuit. A first tape recorder functions as an outer memory of the programmable electronic calculator for storing a program or data processed by the programmable electronic calculator. A second tape recorder functions as a reproduction device for providing audible guidance for the respective steps of the operation, if required."}
{"text": "Most of the automobiles that are in use today have the engine bay located at the front of the vehicle. Typically an internal structural cross-member bridges the front of the bay and supports a radiator that is a part of the engine's cooling system. When the vehicle is equipped with air conditioning, a condenser mounts in front of the radiator. The radiator and condenser are cooled by air that passes through them, the air either being forced through the radiator and condenser by ram air effect when the vehicle is in forward motion, and/or by being drawn through the radiator and condenser by a fan or fans located directly behind the radiator and condenser. As a result, the engine bay is ventilated by air that has been heated by the radiator and condenser. Before this air leaves the engine bay, it is further heated by the heat emitted directly by the engine. Consequently, elevated temperatures can occur in the engine bay and at other locations that are exposed to heated air from the engine bay. These elevated temperatures may be sufficiently high to create thermally induced problems in certain areas or components of an automotive vehicle.\nCommonly assigned, allowed U.S. patent application Ser. No. 07/357,509, filed May 25, 1989, now U.S. Pat. No. 4,979,584, in the name of Herbert N. Charles, discloses a new and unique arrangement for ventilating an engine bay to reduce temperatures. The invention of that patent application involves the use of a ducted fan system to draw cooling air through the radiator and condenser and to convey the hot effluent from the radiator and condenser through a conduit that empties to a location that is outside the engine bay so that the engine bay is not ventilated by the hot effluent. Apertures are provided in the structural cross-member that supports the radiator/condenser heat exchange structure so that the engine bay can be ventilated by ambient ram air when the vehicle is in forward motion. Accordingly, an arrangement of that type provides a means for obtaining significant temperature reductions in the engine bay in comparison to conventional installations where the engine bay is ventilated by hot effluent from radiator/condenser heat exchange structure.\nIt has now been discovered that a ducted fan system like that just described can be used to further enhance the ventilation of the engine bay by the incorporation of an ejector into the ducted fan system at a location downstream of the fan. The effluent from the radiator/condenser heat exchange structure is conducted through the ejector, and the ejector is effective, under certain conditions of operation of the ducted fan, to cause the effluent that passes through the ejector to create a relatively lower pressure zone where the pressure is less than that of a relatively higher pressure zone of the engine bay or of a location adjacent the engine bay that is exposed to heat from the engine bay. The relatively lower pressure zone of the ejector is placed in communication with the relatively higher pressure zone of the engine bay, or adjacent location, for example by means of a ventilating conduit that extends from the relatively higher pressure zone of the engine bay or adjacent location and tees into the ducted fan system at the relatively lower pressure zone of the ejector. As a result, a flow is induced in the ventilating conduit, and that flow draws air from the relatively higher pressure zone of the engine bay or adjacent location into the ejector where the flow entrains with the effluent passing through the ducted fan system to be ultimately discharged from the ducted air system along with the effluent to a location outside the engine bay where the discharge of the effluent is acceptable. It is contemplated that a suitably designed ejector and ventilating conduit can induce in the ventilating conduit flows of as much as about 10% to 15% of the effluent flow without appreciable loss of efficiency in the operation of the ducted fan system. Consequently, the invention now makes it possible to enhance the ventilation of the engine bay, and especially to enhance the ventilation of hot spots in the engine bay or adjacent locations that are exposed to heat from the engine bay.\nCertain components of an automobile that are typically within the engine compartment, fuel lines and batteries for instance, may be located where there is insufficient ventilation to provide acceptable cooling for these components. If such is the case, the present invention can provide a very effective solution which comprises designing an ejector into the ducted fan system, and running a ventilating conduit from the hot spot to the ejector. The solution is especially advantageous because the ejector and the ventilating conduit contain no moving parts. The ejector can be fabricated by conventional plastic molding techniques, and the ventilating conduit can be fabricated in any of a number of conventional ways, such as using flexible or rigid hose or tubing. One of the disclosed embodiments of the invention contemplates an especially efficient use of both materials and available space within the engine bay by integrating the ejector with the ducted fan scroll.\nThe foregoing features, advantages and benefits of the invention, along with additional ones, will appear in the ensuing description and claims which should be considered in conjunction with the accompanying drawings. The drawings disclose a presently preferred embodiment of the invention according to the best mode contemplated at the present time in carrying out the invention."}
{"text": "Market adoption of wireless LAN (WLAN) technology has exploded, as users from a wide range of backgrounds and vertical industries have brought this technology into their homes, offices, and increasingly into the public air space. This inflection point has highlighted not only the limitations of earlier-generation systems, but the changing role WLAN technology now plays in people's work and lifestyles, across the globe. Indeed, WLANs are rapidly changing from convenience networks to business-critical networks. Increasingly users are depending on WLANs to improve the timeliness and productivity of their communications and applications, and in doing so, require greater visibility, security, management, and performance from their network.\nThe rapid proliferation of lightweight, portable computing devices and high-speed WLANs enables users to remain connected to various network resources, while roaming throughout a building or other physical location. The mobility afforded by WLANs has generated a lot of interest in applications and services that are a function of a mobile user's physical location. Examples of such applications include: printing a document on the nearest printer, locating a mobile user or rogue access point, displaying a map of the immediate surroundings, and guiding a user inside a building. The required or desired granularity of location information varies from one application to another. Indeed, the accuracy required by an application that selects the nearest network printer, or locates a rogue access point, often requires the ability to determine in what room a wireless node is located. Accordingly, much effort has been dedicated to improving the accuracy of wireless node location mechanisms.\nThe use of radio signals to estimate the location of a wireless device or node is known. For example, a Global Positioning System (GPS) receiver obtains location information by triangulating its position relative to four satellites that transmit radio signals. The GPS receiver estimates the distance between each satellite based on the time it takes for the radio signals to travel from the satellite to the receiver. Signal propagation time is assessed by determining the time shift required to synchronize the pseudo-random signal transmitted by the satellite and the signal received at the GPS receiver. Although triangulation only requires distance measurements from three points, an additional distance measurement from a fourth satellite is used for error correction.\nThe distance between a wireless transmitter and a receiver can also be estimated based on the strength of the received signal, or more accurately the observed attenuation of the radio signal. Signal attenuation refers to the weakening of a signal over its path of travel due to various factors like terrain, obstructions and environmental conditions. Generally speaking, the magnitude or power of a radio signal weakens as it travels from its source. The attenuation undergone by an electromagnetic wave in transit between a transmitter and a receiver is referred to as path loss. Path loss may be due to many effects such as free-space loss, refraction, reflection, and absorption.\nIn business enterprise environments, most location-tracking systems are based on RF triangulation or RF fingerprinting techniques. RF triangulation calculates a mobile user's location based upon the detected signal strength of nearby access points (APs). It naturally assumes that signal strength is a factor of proximity, which is true a majority of the time. However, the multipath phenomenon encountered in indoor RF environments does present certain difficulties in locating wireless nodes, since reflection and absorption of RF signals affects the correlation between signal strength and proximity. RF fingerprinting compares a mobile station's view of the network infrastructure (i.e., the strength of signals transmitted by infrastructure access points) with a database that contains an RF physical model of the coverage area. This database is typically populated by either an extensive site survey or an RF prediction model of the coverage area. For example, Bahl et al., “A Software System for Locating Mobile Users: Design, Evaluation, and Lessons,” http://research.microsoft.com/˜bahl/Papers/Pdf/radar.pdf, describes an RF location system (the RADAR system) in a WLAN environment, that allows a mobile station to track its own location relative to access points in a WLAN environment.\nThe RADAR system relies on a so-called Radio Map, which is a database of locations in a building and the signal strength of the beacon packets emanating from the access points as observed, or estimated, at those locations. For example, an entry in the Radio Map may look like (x, y, z, ssi (i=1 . . . n)), where (x, y, z) are the physical coordinates of the location where the signal is recorded, and ssi is the signal strength of the beacon signal emanating from the ith access point. According to Bahl et al., Radio Maps may be empirically created based on heuristic evaluations of the signals transmitted by the infrastructure radios at various locations, or mathematically created using a mathematical model of indoor RF signal propagation. To locate the position of the mobile user in real-time, the mobile station measures the signal strength of each of access points within range. It then searches a Radio Map database against the detected signal strengths to find the location with the best match. Bahl et al. also describe averaging the detected signal strength samples, and using a tracking history-based algorithm, to improve the accuracy of the location estimate. Bahl et al. also address fluctuations in RF signal propagation by using multiple Radio Maps and choosing the Radio Map which best reflects the current RF environment. Specifically, one access point detects beacon packets from other access points and consults a radio map to estimate its location, and evaluates the estimated location with the known location. The RADAR system chooses the Radio Map which best characterizes the current RF environment, based on a sliding window average of received signal strengths.\nWhile the RADAR system works for its intended objective, even in this system, location accuracy decreases with the error in detecting the strength of RF signals. For example, individual differences as to how two different wireless nodes detect and report signal strength can cause errors in location, since the Radio Maps assume no error in such measurements. Accordingly, two wireless nodes in the same location that detect different signal strengths will compute different estimated locations. Still further, while the RADAR system allows a mobile station to track its own location, it does not disclose a system that allows the WLAN infrastructure to track the location of wireless nodes, such as rogue access points. Such a system is desirable as it obviates the need for special client software to be installed on the mobile stations.\nThis paradigm shift, however, presents certain problems. As discussed above, the Radio Maps in the RADAR system are constructed from the point of view of a wireless node in an RF environment that includes access points in known locations. In other words, the Radio Maps are constructed based on heuristic and/or mathematical evaluations of the propagation of signals from the access points to a wireless node at a given location. Accordingly, the RADAR system need not assume symmetry of path loss between a given location and the access points in the RADAR system, since the mobile station detects the signal strength of the access points and computes its own location. In addition, since the location of a wireless node is based on path loss, the transmit power of the radio transmitters used to determine location must also be known. In the RADAR system, this is not problematic, since the signals used to determine location are transmitted by access points, whose transmit power can be controlled or easily determined. Estimating location based on signals transmitted by a wireless node, however, can be problematic, since transmit power can vary among wireless device manufacturers, and/or may be individually configured by the mobile user. Furthermore, the more complicated issue of variable transmit power exists not only due to the device or intended application itself, but also due to placement of the wireless node (e.g., a RFID tag, etc.) placed among or inside of other objects, or located within or behind various physical barriers or obstructions.\nOne approach to this problem is to assume symmetry in path loss between a given location in an RF environment and the radio transceivers used to detect signals transmitted by the wireless nodes. Furthermore, these approaches also assume a uniform transmit power for the wireless nodes in light of the fact that legal regulations, as well as current chip set technology, generally places an upper limit on transmit power. These two assumptions, however, can significantly impact the accuracy of locating a wireless node. As discussed above, the RADAR system, for example, finds the location coordinates in the Radio Map that are the best fit based on the detected signal strengths. That is, for each point in the Radio Map, the location metric computes the Euclidean distance between the detected signal strength values and the values in the Radio Map.\nThe following equation provides an illustrative example, assume for didactic purposes that a given wireless node is detected by three access points. The signal strength samples are RSSIap1 RSSIap2, and RSSIap3, while the RF coverage maps for each of the access points are denoted as MAPap1, MAPap2, MAPap3, where the coverage maps include access point signal strength values detected or computed for different locations in a defined region. Again, assuming path loss symmetry and a uniform transmit power, individual error surfaces for each access point can be created based on the signal strength detected at each access point, (RSSIap1, etc.) and the signal strength values in the individual coverage maps (e.g., MAPap1, etc.). That is, the error surface is the difference between the observed signal strength at a given access point less the signal strength values in the coverage map. The locations in this coverage map where the difference is zero are the most likely locations. In many situations, however, the measured signal strengths, RSSIap1 RSSIap2, and RSSIap3, do not match the signal strengths recorded in the coverage maps MAPap1, MAPap2, MAPap3 at any one location. In this case, it is desirable to find the location that is “closest” to matching RSSIap1 RSSIap2, and RSSIap3—in other words, the location that minimizes some function of MAPap1, MAPap2, MAPap3, RSSIap1 RSSIap2, and RSSIap3. Bahl et al., supra, describe several ways in which this function is created, including minimum mean squared error, minimum distance, and minimum Manhattan grid distance. Furthermore, a total error surface, ErrSurf, can be computed based on the sum of the squares (to neutralize positive and negative differences) of the individual error surfaces (i.e., the difference between the detected signal strength values and the signal strength values in each coverage maps), as follows:ErrSurf=[(RSSIap1−MAPap1)^2+(RSSIap2−MAPap2)^2+(RSSIap3−MAPap3)^2]/3In one implementation, the estimated wireless node location is derived from the minimum or minimum of this total error surface.\nHowever, a change in the wireless node's effective transmit power (or, in the RADAR system, inaccuracies in detecting signal strength by the wireless nodes) will adversely affect the accuracy of this metric. For example, an N dB difference between the actual and assumed transmit power of a wireless node would cause a N dB change in the detected signal strengths. Rather than merely shifting the individual signal strength differences for each point in the individual error surfaces up by some fixed amount, the individual differences between the detected signal strengths and the signal strength values in the error surface can change quite dramatically. Indeed, each point in the individual error surfaces are shifted an amount proportional to the dB error. This circumstance moves some areas of the total error surface up relative to others, and some areas of the total error surface down relative to others, significantly altering the shape of the error surface, as well as the location, shape, and size of its minima. It also creates unpredictable error with changes in transmit power. Similar problems will occur for a fixed error in the “link or path loss symmetry” where the path loss from access point to wireless node differs from the path loss from wireless node to access point by some fixed amount due to propagation characteristics, vantage point, and the like. In addition, sources of RF interference typically have unknown transmit powers, and may only partially overlap the frequency band in which wireless nodes operate. Estimating the location of these interference sources requires a method that does not depend entirely on the absolute detected signal strength value.\nIn light of the foregoing, a need in the art exists for a wireless node location mechanism that reduces the errors in computing the location of a wireless node due to commonly occurring circumstances, such as variations in wireless node transmit power, errors in signal strength detection, and/or direction-dependent path loss. Embodiments of the present invention substantially fulfill this need."}
{"text": "In a trade show, an exhibition or a similar environment, in order to help participants to get the most out of the visit, there are usually some electronic tour-guide (purchase guide) systems provided which allow users to browse or hear more information of an object they are interest in through a portable client device. Based on the way these systems decide of which object the user is interest in, generally they can be divided into two types. In the first type of systems, the user should manually key in the identification (ID) of the object he/she is interest in on his/her client device. However, this is quite troublesome for users. The second type of systems is the so-called location-based system, where some kind of location detection mechanism (such as radio frequency identification (RFID)) is used, and the client device will automatically deliver the information of an object to the user when the user is standing approximate to this object. Although this type of system is more convenient than the first one, it has following drawbacks. First, when there are several objects with same or similar distance and directional angle to the user, there is no reliable way to decide which object should be the one the user is interested in. Second, after all, the location approximation has no certain relation with a user's real interest. For example, when the user is facing and looking at a picture hanged on the wall several steps away while there is another exhibited item just next to the user's side, in this case, the object closest to the user is not the one the user is actually interested in. So, the object information delivered to the user by a location-based system is not the information of the object the user is actually interested in."}
{"text": "The inventive concept relates to integrated circuit(s), and more particularly, to computer-implemented method(s) for designing and/or manufacturing integrated circuit(s).\nThe design and manufacture of integrated circuits is a highly complex process. The evolution of integrated circuits is characterized by an increasing integration density of the constituent elements. In order to obtain a reliable integrated circuit at the end of the design and manufacturing process, a number of design constraints (e.g., element proximity, signal timing, power consumption, etc.) are usually defined in order to protect the integrity of the overall design. As integration density increases, the possibility of violating one or more design constraint(s) rises.\nThe evolution of integrated circuits is also characterized by a desire to reduce or minimize the physical size of and/or the power consumed by the integrated circuit (or the incorporating semiconductor chip). Unfortunately, as the number of design constraints for an integrated circuit increases, corresponding efforts to satisfy the design constraints tend to drive up physical size and/or power consumption. For example, efforts to satisfy a particular timing constraint for an integrated circuit may result in the addition of a buffer cell (or analogous delay element) in order to improve hold time margin(s). Alternately or additionally, efforts to satisfy a particular noise constraint for an integrated circuit may result in the addition of a decoupling capacitor. Such efforts tend to increase the physical size and/or power consumption of the integrated circuit."}
{"text": "The present invention relates to compositions and methods for the prevention and treatment of neoplastic and oncogenic disorders, with combinations of certain limonoids, flavonoids and/or tocotrienols. Limonoids are a group of chemically related triterpene derivatives found in the Rutaceae and Meliaceae families. Limonoids are among the bitter principals in citrus juices such as lemon, lime, orange and grapefruit. Flavonoids are polyphenolic compounds that occur unbiquitously in plant foods especially in orange, grapefruit and tangerine. Tocotrienols are present in palm oil and are a form of vitamin E having an unsaturated side chain. In the practice of the cancer prevention and/or treatment of the invention the limonoids, flavonoids and tocotrienols are used to inhibit the development progression and proliferation of cancer cells. Preferred compositions of the invention are those which specifically or preferentially prevent transformation of preneoplastic cells to tumor cells, and prevent or inhibit tumor cell proliferation, invasion and metastasis without general cytotoxic effects.\n2.1 Limonoids\nLimonoids are a group of chemically related triterpene derivatives found in the Rutaceae and Meliaceae families. Limonoids are among the bitter principles found in citrus fruits such as lemons, lime, orange and grapefruit. They are also present as glucose derivatives in mature fruit tissues and seed, and are one of the major secondary metabolites present in Citrus. Limonoids have been found to have anti-carcinogenic activity in laboratory animals. The furan moiety attach to the D-ring is specifically responsible for detoxifying of the chemical carcinogen glutathione-S-transferase enzyme system (Lam, et al., 1994, Food Technology 48:104-108).\nCitrus fruit tissues and byproducts of juice processing such as peels and molasses are sources of limonoid glucosides and citrus seed contain high concentrations of both limonoid aglycones and glucosides. Limonoid aglycones m the fruit tissues gradually disappear during the late stages of fruit growth and maturation.\nThirty-eight limonoid aglycones have been isolated from Citrus. The limonoids are present in three different forms: the dilactone (I) is present as the open D-ring form (monolactone), the limonoate A-ring lactone (II) and the glucoside form (III). Only the monolactones and glucosides are present in fruit tissues. (Hasegawa S. et al., 1994, in Food Phytochemicals for Cancer Prevention I, eds M-T. Huang et al, American Chemical Society, 198-207). ##STR1##\nCompound III is the predominant limonoid glucoside found in all juice samples. In orange juice it comprises 56% of the total limonoid glucosides present, while in grapefruit and lemon juices, it comprises an average of 63% to 66% respectively. Procedures for the extraction and isolation of both aglycones and glucosides have been established to obtain concentrated sources of various limonoids (Lam, L. K. T. et al., 1994, in Food Phytochemicals for Cancer Prevention, eds. M. Huang, T. Osawa, C. Ho and R. T. Rosen, ACS Symposium Series 546, p 209). The use of limonoids in combination with a flavonoid, tocotrienol, a cancer chemotherapeutic agent, or a combination of any one of these agents, has not been reported for the prevention and treatment of neoplastic diseases.\n2.2 Flavonoids.\nEpidemiological studies have shown that flavonoids present in the Mediterranean diet may reduce the risk of death from coronary heart disease (Hertog, M. G. et al., 1993, Lancet: 342, 1007-1011). Soybean isoflavones for example, genistein, which is a minor component of soy protein preparations may have cholesterol-lowering effects (Kurowska, E. M. et al., 1990, J. Nutr. 120:831-836). The flavonoids present in citrus juices such as orange and grapefruit include, but are not limited to, hesperetin and naringenin respectively.\n##STR2## 5 7 3' 4' HESPERETIN OH OH OH OCH.sub.3 NARINGENIN OH OH -- OH\nThe flavonoids preset in tangerine include, but are not limited to tangerevtin or nobiletin. These flavonoids were found to inhibit growth of both estrogen receptor-negative (ER-) and positive (ER+) breast cancer cells in culture and act synergistically with tamoxifen and tocotrienols (Guthrie N. et aL, 1996, Proc. Am. Inst. Cancer Res., Abs.#8).\n ##STR3## 5 6 7 8 4' 5' TANGERETIN O CH.sub.3 O CH.sub.3 O CH.sub.3 O CH.sub.3 O CH.sub.3 -- NOBILETIN O CH.sub.3 O CH.sub.3 O CH.sub.3 O CH.sub.3 O CH.sub.3 O CH.sub.3\n2.3 Tocotrienols in Palm Oil\nTocotrienols are present in palm oil and are a form of vitamin E having an unsaturated side chain. They include, but are not limited to alpha-tocotrienol, gamma-tocotrienol or delta-tocotrienol.\n ##STR4## R1 R2 R3 .alpha.-tocotrienol CH.sub.3 CH.sub.3 CH.sub.3 .gamma.-tocotrienol H CH.sub.3 CH.sub.3 .delta.-tocotrienol H H CH.sub.3\n2.4 Cancer Growth and Chemotherapy\nCancer is a disease of inappropriate tissue accumulation. Chemotherapeutic agents share one characteristic: they are usually more effective in killing or damaging malignant cells than normal cells. However, the fact that they do harm normal cells indicates their potential for toxicity. Animal tumor investigations and human clinical trials have shown that drug combinations produce higher rates of objective response and longer survival than single agents. Combination drug therapy is, therefore, the basis for most chemotherapy employed at present (DeVita, V. T. et. al., 1995, Cancer 35:98).\nCancer treatment requires inhibitions of a variety of factors including tumor cell proliferation, metastatic dissemination of cancer cells to other parts of the body, invasion, tumor-induced neovascularization, and enhancement of host immunological responses and cytotoxicity. Conventional cancer chemotherapeutic agents have often been selected on the basis of their cytotoxicity to tumor cells. However, some anticancer agents have adverse effects on the patient's immune system. Thus it would be greatly advantageous if a cancer therapy or treatment could be developed that would afford non-cytotoxic protection against factors that might lead to progression of tumors.\nBecause hormone therapy as well as chemotherapy is effective in controlling advanced breast cancer, it has been used as an adjuvant to mastectomy in primary breast cancer. Patients with ER+ or ER- tumors benefit from adjuvant chemotherapy. However, tamoxifen used alone as an adjuvant to mastectomy for breast cancer shows benefit in extending disease-free and overall survival (Cummings, F. J. et al., 1985, Ann. Intern. Med. 103;324)."}
{"text": "The invention relates to wireless communications systems. In particular, the invention relates to initialization of a convolutional decoder that has missed a portion of a continuously encoded symbol stream.\nA wireless communication system may comprise multiple remote units and multiple base stations. FIG. 1 exemplifies an embodiment of a terrestrial wireless communication system with three remote units 10A, 10B and 10C and two base stations 12. In FIG. 1, the three remote units are shown as a mobile telephone unit installed in a car 10A, a portable computer remote 10B, and a fixed location unit 10C such as might be found in a wireless local loop or meter reading system. Remote units may be any type of communication unit such as, for example, hand-held personal communication system units, portable data units such as a personal data assistant, or fixed location data units such as meter reading equipment. FIG. 1 shows a forward link 14 from the base station 12 to the remote units 10 and a reverse link 16 from the remote units 10 to the base stations 12.\nCommunication between remote units and base stations, over the wireless channel, can be accomplished using one of a variety of multiple access techniques which facilitate a large number of users in a limited frequency spectrum. These multiple access techniques include time division multiple access (TDMA), frequency division multiple access (FDMA), and code division multiple access (CDMA). An industry standard for CDMA is set forth in the TIA/EIA Interim Standard entitled xe2x80x9cMobile Stationxe2x80x94Base Station Compatibility Standard for Dual-Mode Wideband Spread Spectrum Cellular Systemxe2x80x9d, TIA/EIA/IS-95, and its progeny (collectively referred to here as IS-95), the contents of which are incorporated by reference herein in their entirety. Additional information concerning a CDMA communication system is disclosed in U.S. Pat. No. 4,901,307, entitled SPREAD SPECTRUM MULTIPLE ACCESS COMMUNICATION SYSTEM USING SATELLITE OR TERRESTRIAL REPEATERS, (the \"\"307 patent) assigned to the assignee of the present invention and incorporated in its entirety herein by reference.\nIn the \"\"307 patent, a multiple access technique is disclosed where a large number of mobile telephone system users, each having a transceiver, communicate through base stations using CDMA spread spectrum communication signals. The CDMA modulation techniques disclosed in the \"\"307 patent offer many advantages over other modulation techniques used in wireless communication systems such as TDMA and FDMA. For example, CDMA permits the frequency spectrum to be reused multiple times, thereby permitting an increase in system user capacity. Additionally, use of CDMA techniques permits the special problems of the terrestrial channel to be overcome by mitigation of the adverse effects of multipath, e.g. fading, while also exploiting the advantages thereof.\nIn a wireless communication system, a signal may travel several distinct propagation paths as it propagates between base stations and remote units. The multipath signal generated by the characteristics of the wireless channel presents a challenge to the communication system. One characteristic of a multipath channel is the time spread introduced in a signal that is transmitted through the channel. For example, if an ideal impulse is transmitted over a multipath channel, the received signal appears as a stream of pulses. Another characteristic of the multipath channel is that each path through the channel may cause a different attenuation factor. For example, if an ideal impulse is transmitted over a multipath channel, each pulse of the received stream of pulses generally has a different signal strength than other received pulses. Yet another characteristic of the multipath channel is that each path through the channel may cause a different phase on the signal. For example, if an ideal impulse is transmitted over a multipath channel, each pulse of the received stream of pulses generally has a different phase than other received pulses.\nIn the wireless channel, the multipath is created by reflection of the signal from obstacles in the environment such as, for example, buildings, trees, cars, and people. Accordingly, the wireless channel is generally a time varying multipath channel due to the relative motion of the structures that create the multipath. For example, if an ideal impulse is transmitted over the time varying multipath channel, the received stream of pulses changes in time delay, attenuation, and phase as a function of the time that the ideal impulse is transmitted.\nThe multipath characteristics of a channel can affect the signal received by the remote unit and result in, among other things, fading of the signal. Fading is the result of the phasing characteristics of the multipath channel. A fade occurs when multipath vectors add destructively, yielding a received signal that is smaller in amplitude than either individual vector. For example if a sine wave is transmitted through a multipath channel having two paths where the first path has an attenuation factor of X dB, a time delay of xcex4 with a phase shift of \"THgr\" radians, and the second path has an attenuation factor of X dB, a time delay of xcex4 with a phase shift of \"THgr\"+xcfx80 radians, no signal is received at the output of the channel because the two signals, being equal amplitude and opposite phase, cancel each other. Thus, fading may have a severe negative effect on the performance of a wireless communication system.\nTypically, modern communication systems use coding to improve immunity to interference and wireless channel noise. Additionally, coding may increase system capacity and improve security. Generally, an information signal is first converted into a form suitable for efficient transmission over the wireless channel. Conversion or modulation of the information signal involves varying a parameter of a carrier wave on the basis of the information signal in such a way that the spectrum of the resulting modulated carrier is confined within the channel bandwidth. At a remote unit, the original message signal is replicated from a version of the modulated carrier received following propagation over the wireless channel. Such replication is generally achieved by using an inverse of the modulation process employed by the base station.\nThe field of data communications is particularly concerned with optimizing data throughput of a transmission system with a limited signal to noise ratio (SNR). The use of error correcting circuitry, such as encoders and decoders, allows system tradeoffs to be made. For example, smaller SNRs or higher data rates may be used with a particular wireless channel which maintains the same bit error rate (BER).\nOne class of encoders is known as a convolutional encoder. As is well known in the art, a convolutional encoder converts a sequence of input data bits to a codeword based on a convolution of the input sequence with itself or with another signal. Convolutional encoding of data combined with a convolutional decoder is a well known technique for providing error correction coding and decoding of data. One type of convolutional decoder typically used is a Viterbi decoder.\nCoding rate, constraint length, and generating polynomials are used to define a convolutional decoder. A coding rate (k/n) corresponds to the number of coding symbols produced (n) for a given number of input bits (k). A constraint length (K) is defined as the length of a shift register used in convolutional encoding of data. Convolutional codes add correlation to an input data sequence by using delay elements (i.e., shift registers) and modulo adders. Taps between the delay elements may terminate at modulo adders forming a desired generating polynomial.\nFIG. 2 is a block diagram of a convolutional encoder 20. The encoder 20 shown contains a shift register 22 tapped at various positions 23A through 23N. The shift register taps terminate at one or more of the modulo-2 adders 24 and 25 forming generator functions g0 and g1. Different generating polynomials can be formed by the selection of which taps terminate at the modulo-2 adders.\nBits enter the encoder at its input 26 one at a time. The outputs of the generator functions are the encoded output symbols C0 and C1. Each of the two generator functions g0 and g1 produces an output symbol for each input bit, which corresponds to a code rate of xc2xd. For an encoder with three generator functions, the code rate is ⅓ and a code rate of {fraction (1/n )} requires n generator functions. \nA coding rate of xc2xd has become one of the most popular rates, although other code rates may be used. A constraint length of nine (K=9) is typical in convolutional coding schemes. The convolutional encoder can be thought of as a Finite Impulse Response filter with binary coefficients and length Kxe2x88x921. This filter produces a symbol stream with 2Kxe2x88x921 possible states.\nA basic principle of the Viterbi algorithm is to take a convolutionally encoded data stream that has been transmitted over a noisy wireless channel and use a finite state machine to efficiently determine the most likely sequence that was transmitted. A K=9 Viterbi decoder can be thought of as a machine that hypothesizes which of each 256 (2(Kxe2x88x921)) possible states the encoder could have been in given the symbols received. The probability that the encoder transitioned from each of those states to the next set of 256 possible encoder states is determined. The probabilities are represented by quantities called metrics, which are proportional to the negative of the logarithm of the probability. The sum of the metrics is therefore equivalent to the reciprocal of the product of the probabilities. Thus, smaller metrics correspond to higher probability events.\nThere are two types of metrics: state metrics, sometimes called path metrics; and branch metrics. The state metric represents the probability that the received set of symbols leads to the state with which it is associated. The branch metric represents the conditional probability that the transition from one state to another occurred assuming that the starting state was actually the correct state and given the symbol that was actually received.\nIn a rate {fraction (1/N )} encoder, there are two possible states leading to any other state, each corresponding to the occurrence of a zero or a one in the right-most bit of the convolutional encoder shift register. The decoder decides which is the more likely state by an add-compare-select (ACS) operation. Add refers to adding each state metric at the preceding level to the two branch metrics of the branches for the allowable transitions. Compare refers to comparing the pair of such metric sums for paths entering a state (node) at the given level. Select refers to selecting the lesser of the two and discarding the other. Thus, only the winning branch, i.e., the branch with the highest probability (smallest metric), is preserved at each node, along with the node state metric. If the two quantities being compared are equal, either branch may be selected, for the probability of erroneous selection will be the same in either case. \nThe Viterbi algorithm is a computationally efficient method of updating the conditional probabilities of the best state and the most probable bit sequence transmitted from the possible 2Kxe2x88x921 states. In order to compute this probability, all 2Kxe2x88x921 states for each bit must be computed. The resulting decision bits from each of these computations is stored as a single bit in a path memory.\nA chain-back operation, an inverse of the encoding operation, is performed in which the C decision bits are used to select an output bit, where C is the chainback distance. After many branches the most probable path will be selected with a high degree of certainty. The path memory depth must be sufficiently long to be governed by the signal-to-noise ratio and not the length of the chain-back memory.\nThough it is not necessary for analyzing either the code characteristics or the performance of the optimal decoder, it is useful in understanding both to exhibit the code on a trellis diagram. The term xe2x80x9ctrellisxe2x80x9d is a term which describes a tree in which a branch not only bifurcates into two or more branches but also in which two or more branches can merge into one. A trellis diagram is an infinite replication of the state diagram for an encoder. The nodes (states) at one level in the trellis are reached from the node states of the previous level by the transition through one branch, corresponding to one input bit, as determined by the state diagram. Any codeword of a convolutional code corresponds to the symbols along a path (consisting of successive branches) in the trellis diagram.\nA simple embodiment of the encoder of FIG. 2 is illustrated in FIG. 3. FIG. 3 illustrates a convolutional encoder with a code rate of xc2xd and a constraint length of 3. As shown in FIG. 3, the convolutional encoder has three taps 31, 32 and 33. The taps terminate at two modulo 2 adders 35 and 36 forming generator functions g0=510 and g1=710. The output of the generator functions become the encoded output symbols C0 and C1, respectively.\nFIG. 4 is a trellis diagram showing the possible paths of the convolutional encoder illustrated in FIG. 3. The encoder is assumed to begin in the zero state. Each possible state is represented in the trellis diagram by a node 42. In each state the next input bit into the encoder may be either a zero or a one and a corresponding set of symbols are generated in each generator function. In FIG. 4, input bits 44 at each state are represented on their associated path. The output code symbols C0, indicated as 46, and C1, indicated as 47, generated from the input of each bit are represented in the diagram on the associated path. As illustrated in this simple example, each set of code symbols received at the remote unit is influenced from previously input data bits at the encoder. Thus, in typical operation, a convolutional decoder receives a continuous uninterrupted stream of code symbols with each symbol influenced by the preceding input data.\nIn a typical CDMA communication system, remote units only sporadically establish bidirectional communication with a base station. For example, a cellular telephone remains idle for significant periods of time when no call is in process. To ensure that any message directed to a remote unit is received, the remote unit must continuously monitor the communication channel even while it is idle. For example, while idle, the remote unit monitors the forward link channel from the base station to detect incoming calls. During such idle periods, the cellular telephone continues to consume power to sustain the elements necessary to monitor for signals from the base stations. Many remote units are portable and are powered by an internal battery. For example, personal communication system (PCS) handsets are almost exclusively battery-powered. The consumption of battery resources by the remote unit in idle mode decreases the battery resources available to the remote unit when a call is placed or received. Therefore, it is desirable to minimize power consumption in a remote unit in the idle state and thereby increase battery life.\nOne means of reducing remote unit power consumption in a communication system is disclosed in U.S. Pat. No. 5,392,287, entitled APPARATUS AND METHOD FOR REDUCING POWER CONSUMPTION IN A MOBILE COMMUNICATION RECEIVER (the \"\"287 patent), assigned to the assignee of the present invention and hereby incorporated in its entirety herein by reference. In the \"\"287 patent, a technique for reducing power consumption in a remote unit operating in an idle mode (i.e. a remote unit which is not engaged in bidirectional communication with a base station) is disclosed. In idle, each remote unit periodically enters an xe2x80x9cactivexe2x80x9d state during which it prepares to and receives messages on a forward link communication channel. In the time period between successive active states, the remote unit enters an xe2x80x9cinactivexe2x80x9d state. During the remote unit\"\"s inactive state, the base station does not send any messages to that remote unit, although it may send messages to other remote units in the system that are in the active state.\nAs disclosed in the \"\"287 patent, a base station broadcast messages which are received by all remote units within the base station coverage area on a xe2x80x9cpaging channel.xe2x80x9d All idle remote units within the base station coverage area monitor the paging channel. The paging channel is divided in the time dimension into a continuous stream of xe2x80x9cslots.xe2x80x9d Each remote unit operating in slotted mode monitors only specific slots which have been assigned to it as active (assigned) slots. The paging channel continually transmits convolutional encoded messages in numbered slots, repeating the slot sequence, such as for example, every 640 slots. When a remote unit enters the coverage area of a base station, or if a remote unit is initially powered on, it communicates its presence to a preferred base station. Typically the preferred base station is the base station which has the strongest pilot signal as measured by the remote unit.\nThe preferred base station, along with a plurality of geographically near neighboring base stations, assign a slot, or a plurality of slots, within their respective paging channels, for the remote unit to monitor. The base station uses the slots in the paging channel to transmit control information to a remote unit, if necessary. The remote unit may also monitor a timing signal from the preferred base station allowing the remote unit to align, in the time dimension, to the base station slot timing. By aligning in the time dimension to the preferred base station slot timing, the remote unit can determine when a paging channel slot sequence begins. Thus, knowing when the paging channel slot sequence begins, which slots are assigned for it to monitor, the total number of slots in the repetitive paging channel sequence of slots, and the period of each slot, the remote unit is able to determine when its assigned slots occur.\nGenerally, the remote unit is in the inactive state while the base station is transmitting on the paging channel in slots which are not within the remote unit\"\"s assigned set. While in the inactive state, the remote unit does not monitor timing signals transmitted by the base station, maintaining slot timing using an internal clock source. Additionally, while in the inactive state the remote unit may remove power and/or clocks from selected circuitry, such as, for example, circuits which monitor the wireless channel and the decoder. Using its internal timing, the remote unit transits to its active state a short period of time before the next occurrence of an assigned slot.\nWhen transiting to the active state, the remote unit applies power to circuitry that monitors the wireless channel. After the remote unit has reacquired the base station, the remote unit begins receiving a stream of coded symbols and clocks the coded symbols into the decoder. The decoder uses the coded symbols to continue building a trellis stored in the decoder. However, because the stream of coded symbols has been interrupted, the symbol codes being received by the remote unit have no relationship to the symbols that built the trellis stored within the decoder. Therefore the remote unit must receive sufficient code symbols prior to its assigned slot to ensure that proper decoding of the code word is accomplished. For example, the paging channel used in IS-95 is continuously encoded with a constraint length 9 convolutional code. A decoder used to decode the IS-95 paging channel may need to decode 116 data bits to properly initialize its state metrics and insure valid data is output from the decoder.\nWhen the remote unit enters the active state, it may receive messages in its assigned slots in the paging channel and respond to commands from the base station. For example, the remote unit may be commanded to activate a xe2x80x9ctrafficxe2x80x9d channel to establish a bi-directional communication link for conducting subsequent voice communication in response to an incoming call. If there is no message from the base station, or no command requesting the remote unit to remain active, at the end of the assigned slot the remote unit returns to the inactive state. In addition, the remote unit returns to the inactive state immediately if commanded to do so by the base station.\nTherefore, there is a need in the art for a method and apparatus to decrease the number of code symbols required to properly decode an interrupted stream of code words.\nThe invention addresses these and other needs by providing a system and method wherein the convergence of a convolutional decoder is improved. In one aspect of the invention, the remote unit comprising the convolution decoder receives an interrupted stream of code symbols. Prior to decoding the symbols, the state metrics of the trellis residing within the decoder are initialized.\nIn one aspect of the invention, when the remote unit receives the interrupted stream of code symbols, the pattern of code symbols transmitted just prior to the code symbols received are known. The state metrics of the trellis are then biased towards states which would have been valid only if the previous code symbols would have been received. The remaining invalid states of the trellis are initialized with a high state metric. Therefore, during the decoding process, the decoder is biased toward the valid states.\nIn another aspect of the invention, the metrics of all states in the trellis residing in the encoder are initialized to 0 or some other constant value. Therefore, there is not a bias towards any particular state within the trellis."}
{"text": "The grid of electric power distribution in the United States includes main, secondary and tertiary high voltage (high tension) cables installed underground in conduits and ducts, normally accessible by manholes. In a large Metropolitan area, such as New York City, servicing millions of customers, there are numerous cables installed in cable bundles in a particular duct or conduit. Under normal operation, the cable jacket and insulation provides adequate thermal and electrical isolation to prevent arcing (electrical discharge) between cables.\nHigh ambient temperatures and low air circulation in manholes, especially during the summer months, combined with the heat generated by the electric cables, can cause electrical breakdown of cable insulation, arcing and cable fires that rapidly transfer fire and heat to adjacent cables, causing further arcing. The result is a chain reaction that causes multiple cable failures interrupting power to all, or nearly all, cables in a manhole, duct or conduit. Damage spread over a wide area results in power outages to many people, requiring days or weeks to repair and restore service.\nAdditionally, the high levels of smoke and toxic products of combustion will endanger the lives of many people in adjacent areas, such as subway systems, as well as many individuals attempting to enter the affected areas to service or replace damaged cables.\nCurrent cable wraps provide sufficient high voltage protection under normal conditions, but little, if any, fire and heat protection. Many of these wraps were tested and certified more than thirty years ago based on very low standards of fire protection and no requirements for limitations on smoke or toxic products of combustion. The more recently developed wraps continue to use standards for fire protection that do not adequately model the severe conditions in installations today. Further, they do not provide fire, high voltage and environmental resistance necessary for adequate protection in severe conditions.\nThere is a need for high performance, intumescent and high voltage wrap to isolate each cable under ordinary conditions and more importantly to provide the necessary isolation for cables in the event of overload, dielectric breakdown, and power surge in intense heat that could otherwise lead to a fire. The wrap isolates a cable failure and power can be switched, either by an operator or automatically by the system, to maintain continuity. Customers serviced by the failed cable, or adjacent cables, experience no disruption of service or only momentary disruption due to a switch over.\nIntumescent coatings have been known for many years and have been used to provide thermal protection for many substrates including; wood and wood products, metals, fiberglass, and many types of plastics. However, an intumescent coating applied to these aforementioned substances is impractical for the fire and high voltage protection of electrical cables. Currently there is no single fire and high voltage wrap for electrical cables having all the following features:\n1) Providing an effective cable wrap that substantially reduces, or eliminates fire spread along cables and heat transmission to adjacent cables.\n2) Provide an effective intumescent fire-retardant coating that substantially reduces, or eliminates, smoke and toxic products of combustion from cable insulation when exposed directly or indirectly to a fire.\n3) Provide an effective intumescent fire-retardant coating that is noncombustible and will prevent or eliminate ignition of cable insulation.\n4) Provide a cable wrap with superior mechanical properties, including flexibility that allows easy handling, and wrapping of electrical cables of varying diameters.\n5) Provide a durable and resistant coating to resist abrasion, impact, water, hydrocarbons, chemicals and other environmental factors associated with underground cable installation.\n6) Provide a wrap that has high dielectric breakdown strength providing electrical insulation to 50 kVolts.\n7) Provide a wrap that has very low ampacity deration, to maintain the high electric current carrying capacity of the cables.\n8) Provide a cable wrap that is easily manufactured for mass production."}
{"text": "The present invention relates to a nonvolatile memory device, a method of manufacturing the device, and a method of driving the device, and more particularly, to a nonvolatile memory in which the coupling ratio of the memory cells is increased without increasing cell size, through the structure and operation of a \"program assist plate,\" thereby lowering the operating voltage and increasing the operating speed of the device. The invention may be used in many different types of nonvolatile memory devices, including NAND, NOR, AND, DINOR and other devices.\nIn a NOR-type electrically erasable programmable read-only memory(EEPROM), two facing memory cells share one bitline contact and one source line, and the memory cells in a row are connected to one bitline. Thus, it is difficult to highly integrate the NOR-type structure, although its high cell current allows it to operate at high speeds.\nIn a NAND-type structure, two cell strings share one bitline contact and one source line. In one cell string, a plurality of cell transistors are connected in series to the bitline. Accordingly, the NAND-type structure can easily obtain a high level of integration, but it is typically slower than the NOR-type structure due to its low cell current. Because the NAND-type memory cell can be more highly integrated than the NOR-type memory cell, it is generally preferable to employ the NAND-type memory structure for increasing the capacity of a memory device. However, this invention is not limited to application in only NAND type devices\nThe EEPROM NAND string structure and the basic operation of the NAND-type EEPROM are described below, referring to the accompanying drawings.\nFIG. 1 is a plan view showing the layout with respect to one string in a typical NAND-type nonvolatile memory device, and FIG. 2 is an equivalent circuit diagram of the structure shown in FIG. 1.\nReferring to FIGS. 1 and 2, each string of a NAND-type nonvolatile memory device is formed by sequentially connecting a string selection transistor S1, a plurality of cell transistors C1, . . . , Cn and a source selection transistor S2 in series between the bitline B/L and a source line S/L in an area represented by a width x and a length y.\nFIG. 3A is a plan view of a transistor cell used in forming each string of the nonvolatile memory device, and FIG. 3B is a sectional view taken along line I-I' of FIG. 3A.\nIn FIG. 3A, reference numeral 26 indicates a mask pattern for forming an active region, reference numeral 24 indicates a mask pattern for forming a control gate, and reference numeral 22 indicates a mask pattern for forming a floating gate.\nReferring to FIG. 3B, each transistor cell C1, . . . , Cn of FIG. 1 in the string consists of a floating gate 32, a control gate 34 and a N-type source/drain 36, which are sequentially deposited on a P-type semiconductor substrate 30, with an interdielectric layer inserted therebetween. The programming, erasing and reading of a NAND-type nonvolatile memory device having this structure is described below.\nThe NAND-type nonvolatile memory is programmed by tunneling an electric charge from a channel region of the cell transistor to the floating gate thereof, to thereby store information. For example, if information is to be programmed or stored in the first transistor cell C1, power supply voltage Vcc is applied to the gate of string select transistor S1, thereby turning on string select transistor S1, and 0V is applied to the gate of source select transistor S2, thereby turning off source select transistor. With reference to FIG. 3B, a programming voltage Vpgm is applied to the control gate 34 of the first transistor cell C1, to thereby generate tunneling. Accordingly, an electric charge in the channel region of the substrate 30 moves to the floating gate 32, to thereby change the threshold voltage Vth of the first transistor cell C1.\nAfter programming, transistor cell C1 will have (approximately) one of two different threshold voltages depending on the charge transferred to the floating gate 32. The first and second threshold voltages may correspond to either a \"1\" or \"0\" in a two-state memory device. In a multi-state memory device more than two threshold voltages may be used, thereby storing more than one bit per cell.\nA read operation is used to determine the programmed state of the NAND memory cell. For example, referring to FIG. 2, when reading information stored in the first cell transistor C1, the bitline B/L is precharged with a predetermined voltage between approximately 1 V.about.Vcc. Then, Vcc is applied to each control gate of the string select transistor S1, the source select transistor S2 and unselected cell transistors C2, . . . , Cn, (i.e., each cell transistor except for C1) to thereby turn-on the transistors. Approximately 0 V is applied to the control gate of the selected first cell transistor C1, which is between a first threshold voltage of approximately -3 V when a \"1\" is stored in the cell and a second threshold voltage of approximately 1V when a \"0\" is stored in cell C1. Thus, if the first cell transistor C1 is turned on, and a current is sensed between the bitline B/L and source line S/L, the state of the first cell transistor C1 is determined as \"1\". However, if the first cell transistor C1 is turned off, and no (or very little) current is sensed between the bitline B/L and the source line S/L, the state of the first cell transistor C1 is determined as \"0\". Alternatively, no current could correspond to a \"1\" and a sensed current could correspond to a \"0\".\nThe erasing operation is performed by tunneling an electric charge from the floating gate 32 to the channel region of the substrate 30 (FIG. 3B), thereby erasing information stored in the cell. For example, referring to FIG. 2, when information is to be erased from the first cell transistor C1, the cell string is placed in a floating state by disconnecting it from the bitline B/L and the source line S/L by turning off the string select transistor S1 and source select transistor S2. A voltage of 0V is applied to all 25 wordlines of a selected block of memory cells C1, C2, . . . Cn. Further, referring to FIG. 3B, an erase voltage Verase is applied to the substrate 30, thereby generating tunneling from the floating gate 32 to the substrate 30. Thus, the electric charge on the floating gate 32 is moved to the substrate 30, thereby changing the threshold voltage of the selected memory cells.\nIn the operation of the nonvolatile memory device described above, a high-voltage of approximately 20V is required to program or erase the memory cells by Fowler-Nordheim (referred to as \"F-N\") tunneling. A charge pumping circuit is required to supply a high voltage for programming and erasing, which results in increased chip size and power consumption. Accordingly, in order to increase the density of a nonvolatile memory device, it is important to increase the efficiency of both erasing and programming, and thereby lower the power requirements for Vpgm and Verase.\nIn order to enhance the operating characteristics without lowering the reliability of the nonvolatile memory device, the capacitance of the structure corresponding to the interdielectric layer deposited between the control gate 34 and the floating gate 32 must be increased, and the program/erase voltage must be lowered. The capacitance may be increased by reducing the thickness of the interdielectric layer or increasing the contact area of the control gate 34 and the floating gate 32. If the capacitance is increased by reducing the thickness of the interdielectric layer, the data retention capability of the nonvolatile memory device is reduced, and the insulation of the interdielectric layer may be broken during programming and erasing. In addition, the process for producing an interdielectric layer of reduced thickness is difficult. However, a method has recently been developed for increasing the contact area between the control gate 34 and the floating gate 32.\nFIG. 4 is a plan view showing a layout of a conventional NAND-type nonvolatile memory device, disclosed in IEDM Tech. Dig. 1994, pp. 61-64, which is incorporated by reference herein. This article discloses a structure and method for obtaining high-integration and increased capacitance with respect to the interdielectric layer by increasing the effective surface area.\nIn FIG. 4, reference numeral 40 denotes a mask pattern for defining an active region, reference numeral 42 denotes a mask pattern for forming a floating gate, reference numeral 44 denotes a mask pattern for forming a control gate, and reference numeral 46 denotes a mask pattern for forming a bitline contact. The mask pattern 42 for forming the floating gate completely overlaps with the mask pattern 40 for defining an active region. That is, the floating gate is self-aligned on the active region, which leads to high-integration.\nFIG. 5 is a sectional view taken along line II-II' of FIG. 4, where reference numeral 50 denotes a semiconductor substrate, reference numeral 52 denotes a floating gate, reference numeral 54 denotes an interdielectric layer, reference numeral 56 denotes a control gate, and reference numeral 58 denotes an isolation film. In this nonvolatile memory device, the area of the interdielectric layer 54 between the floating gate 52 and the control gate 56 is determined only by the width of the active region of the cell transistor. The active regions are the portions of the substrate 50 located between adjacent isolation films 58. Note that the floating gate 52 is not formed on the isolation film 58. Accordingly, the capacitance related to the interdielectric layer 54 is lowered, and the resulting device requires a high voltage for programming and erasing.\nTo solve the above problems, the thickness of the floating gate 52 of FIG. 5 is increased. However, the thick floating gate 52 structure causes two problems. First, when the control gate 56, the interdielectric layer 54, and the floating gate 52 are patterned in accordance with this process, a vertically high interdielectric layer 54 must be formed on the sidewalls of the thick floating gate 52. In addition, it is difficult to etch the thick floating gate 52 in the source/drain region (not shown) of a cell transistor."}
{"text": "The plasma display technology is relatively new and has generated a whole new set of requirements for regulated power supplies. A plasma display will typically require a power supply which generates highly regulated 200 volt D.C. and 5 volt D.C. outputs. The supply voltages are generally isolated from other voltages within the computer and the power supply must generally be current limited and short circuit protected. An additional requirement of the power supply used with a plasma display is that the 200 volt D.C. output must be enabled only after other computer operational voltages have stabilized. If other operational voltages fail for any reason, the 200 volt D.C. voltage for the plasma display must be disabled. This sequencing of the 200 volt D.C. output is required to prevent damage to the plasma display.\nAn adjustable three terminal regulator, due to its not having a ground end, can regulate high voltages as long as its differential voltage rating is not exceeded. In Linear Technology Data Book, 1986, at pages 71-72, a circuit is described which will protect the regulator from a short circuit condition but, as stated in the reference, such a scheme for high voltage regulation is limited by the power dissipation capabilities of the device in series with the regulator.\nWith an input voltage of 250 volts and a typical current of 200 milliamps, the power dissipated is approximately 50 watts. Should a short circuit condition exist for an extended period, which may occasionally occur, an excessive amount of heat is generated requiring an unacceptably large power transistor and heat sink.\nAn attempt to reduce the power dissipation by causing the current limit circuitry to oscillate causes the load change to be reflected to other outputs, causing high ripple content at a frequency approximating the current limit oscillations.\nFurther, providing an input voltage to the regulator that will not exceed the safe operating levels of the regulator becomes a concern due to the voltage overshoot caused by transformer leakage inductance. Left unimpeded, this inductance can allow the input voltage to vary by as much as 30%.\nThe plasma display technology further requires a sequencing or timing control of the 200 volt D.C. output to the display to prevent possible damage o the display. The display voltages must be enabled only after all other computer operational voltages have been established and have stabilized at their normal level. In the event one of those operational voltages fail, the display voltages must be disabled. This sequencing or timing requirement requires coordination of the display voltages with other operational voltages within the computer."}
{"text": "From DE 10 2007 056 638 A1 a generic device is known for assembling a composite arrangement, consisting of a plurality of functional elements having an aperture for a shaft, in particular cams, balancing masses, toothed wheels and/or bearings, in a predetermined angular position on the shaft, wherein the device has a plurality of retaining devices intended respectively for a functional element.\nFrom DE 10 2008 064 194 A1 an apparatus is known for the positioning of a plurality of functional elements having an aperture for a shaft, in particular cams, in a predetermined angular position on the shaft, wherein the apparatus has a plurality of mounts intended respectively for a functional element, which mounts are equipped respectively with a moulding fixing the angular position of the respective functional element corresponding to the angular position on the shaft. The mounts are able to be positioned here such that the apertures of the functional elements lie on a shared horizontal line.\nFrom DE 10 2009 060 350 A1 an apparatus is known for assembling a shaft carrying functional elements, wherein the apparatus comprises a machine platform, on which a plurality of positioning discs for the aligned, correct positioning of the functional elements is arranged in such a manner that a shaft can be pushed in. The positioning discs are reversibly fixed on a frame, which in turn is reversibly fixed on the machine platform. This is intended to make possible rapid changing of a production process by having several frames available."}
{"text": "The present invention relates to a pendulum mechanism, adapted to be attached to the casing of a clockwork for use together with the clockwork, which operates independently of the clockwork, and which comprises a pivot bearing, a pendulum arm mounted on the bearing for swinging motion, and an electrical drive mechanism for the pendulum comprising a permanent magnet carried by the pendulum adjacent its lower end and cooperating with a coil mounted on a circuit board located below the clockwork casing and carrying the electronics for driving the pendulum.\nAn autonomously operating pendulum arrangement, which can be attached to a crystal-controlled clockwork, is discussed in Haag et al U.S. Pat. No. 4,043,118, assigned to the assignee of the instant application. This known arrangement consists of a pendulum housing adapted to be fastened to the lower side of the clockwork casing and carrying a circuit board on which are mounted the electrical components required to drive the pendulum. The arrangement includes a bearing device for the pendulum arm which bearing device is attached as a separate element to the upper side of the clockwork casing, and a device for suspending the pendulum is provided at the lower side of the pendulum arm. The housing is also provided with a lateral arm on which the bell of a clock-striking system can be mounted.\nThis known arrangement has the disadvantage that the proper functioning of its various components, i.e., the overall pendulum arrangement as well as the associated chime system, can be checked out for proper operation only when all of the separate parts of the pendulum have been assembled on the clockwork casing. There is, in addition, the disadvantage that the pendulum arm is substantially U-shaped and envelopes the clockwork laterally, and this requires a comparatively large amount of space for the swinging motion of the pendulum. The known arrangement, moreover, cannot be used in conjunction with comparatively small pendulum clocks, for example table clocks; and the lateral placement of the bell of the clock-striking system precludes use of the arrangement in small-sized classic clocks.\nThe present invention is intended to obviate these disadvantages of the prior art, and is concerned with a novel autonomous pendulum mechanism which can be connected to the casing of the clockwork in the form of a completely separate assembled structural unit whose functioning can be checked out independently of the clockwork. Moreover, the unitary pendulum mechanism of the present invention is so arranged that a chiming system can be incorporated therein as part of the pendulum unit, so that its proper operation can also be checked out independently of the clockwork."}
{"text": "In a communication system, a light modulator may amplitude modulate a light source based on a data signal, and the resulting modulated light, which conveys the data signal, is transmitted to a distant light receiver. In some situations it is advantageous to be able to transmit multiple data signals with different data rates, simultaneously. A conventional light modulator is limited in that it is not capable of simultaneously amplitude modulating a light source with such data signals having the different data rates. Instead, multiple conventional light modulators must be operated in parallel each to transmit a different one of the data signals. This wastes power and increases costs.\nIn the drawings, the leftmost digit(s) of a reference number identifies the drawing in which the reference number first appears."}
{"text": "The allenyl .beta.-lactam compound represented by the general formula (4) given below is conventionally prepared, for example, by reacting a tertiary organic base with the starting material in an organic solvent according to the process disclosed in JP-A-282359/1992. However, this compound is unstable when present in the resulting reaction mixture, which therefore usually requires cumbersome repeated procedures for extraction and concentration after the completion of the reaction. These procedures require time in the case of a large scale production, consequently entailing problems such as a marked reduction in the yield of the isolated product. Thus, processes still remain to be developed which are satisfactorily feasible for the preparation of allenyl .beta.-lactam compounds.\nReports have been made on widely acceptable processes for preparing 3-halogenated cephem derivatives represented by the general formula (6) given below. These processes include a process which uses a compound of the general formula (7), i.e., 3-hydroxycephem compound, as the starting material and involves the conversion of hydroxyl group to trifluoromesyloxy group first and the subsequent reaction with a lithium halide as described in J. Org. Chem., 54, 4962(1989), a process wherein a reactive chlorine or bromine compound (such as phosphorus trichloride, phosphorus oxychloride or thionyl bromide) is reacted with a 3-- hydroxycephem compound in dimethylformamide as disclosed in JP-A-116095/1974), and further a process wherein an alkali metal salt or alkaline-earth metal salt of a halogen is reacted with an allenyl .beta.-lactam compound as disclosed in JP-A-282387/1992 ##STR1## wherein R.sup.1, R.sup.2 and R.sup.3 are as defined below.\nThe first of the processes requires the use of the 3-hydroxycephem compound as the starting material which compound itself is difficult to prepare, so that the process is in no way practically feasible. The second process inevitably forms 3-sulfonylcephem or 3-thiocephem as a by-product due to the recombination of sulfinate ion or thiolate ion which is released on ring closure, consequently giving the desired 3-halogenated cephem derivative in a yield of as low as up to 70%.\nAn object of the present invention is to provide a process wherein a .beta.-lactam compound represented by the general formula (1) is used as the starting material for preparing an allenyl .beta.-lactam compound of the general formula (4) as isolated with stability and high purity in a high yield by a simplified procedure.\nAnother object of the invention is to overcome the drawbacks of the conventional production processes described and to provide a widely-useful process for preparing the desired 3-halogenated cephem derivative in a high yield with a high purity."}
{"text": "The use of computing devices, such as cellular phones and personal digital assistants (PDAs) has grown rapidly. Such devices provide many different functions to users through different types of interfaces, such as keypads and displays. Some computing devices utilize motion as an interface by detecting tilt of the device by a user. Some implementations of a motion interface involve tethering a computing device with fishing lines or carrying large magnetic tracking units that require large amounts of power."}
{"text": "Certain embodiments of the present invention are directed to integrated circuits. More particularly, some embodiments of the invention provide a system and method for stage-based control related to TRIAC dimmer. Merely by way of example, some embodiments of the invention have been applied to driving one or more light emitting diodes (LEDs). But it would be recognized that the invention has a much broader range of applicability.\nA conventional lighting system may include or may not include a TRIAC dimmer that is a dimmer including a Triode for Alternating Current (TRIAC). For example, the TRIAC dimmer is either a leading-edge TRIAC dimmer or a trailing-edge TRIAC dimmer. Often, the leading-edge TRIAC dimmer and the trailing-edge TRIAC dimmer are configured to receive an alternating-current (AC) input voltage, process the AC input voltage by clipping part of the waveform of the AC input voltage, and generate a voltage that is then received by a rectifier (e.g., a full wave rectifying bridge) in order to generate a rectified output voltage.\nFIG. 1 shows certain conventional timing diagrams for a leading-edge TRIAC dimmer and a trailing-edge TRIAC dimmer. The waveforms 110, 120, and 130 are merely examples. Each of the waveforms 110, 120, and 130 represents a rectified output voltage as a function of time that is generated by a rectifier. For the waveform 110, the rectifier receives an AC input voltage without any processing by a TRIAC dimmer. For the waveform 120, an AC input voltage is received by a leading-edge TRIAC dimmer, and the voltage generated by the leading-edge TRIAC dimmer is received by the rectifier, which then generates the rectified output voltage. For the waveform 130, an AC input voltage is received by a trailing-edge TRIAC dimmer, and the voltage generated by the trailing-edge TRIAC dimmer is received by the rectifier, which then generates the rectified output voltage.\nAs shown by the waveform 110, each cycle of the rectified output voltage has, for example, a phase angel (e.g., ϕ) that changes from 0° to 180° and then from 180° to 360°. As shown by the waveform 120, the leading-edge TRIAC dimmer usually processes the AC input voltage by clipping part of the waveform that corresponds to the phase angel starting at 0° or starting at 180°. As shown by the waveform 130, the trailing-edge TRIAC dimmer often processes the AC input voltage by clipping part of the waveform that corresponds to the phase angel ending at 180° or ending at 360°.\nVarious conventional technologies have been used to detect whether or not a TRIAC dimmer has been included in a lighting system, and if a TRIAC dimmer is detected to be included in the lighting system, whether the TRIAC dimmer is a leading-edge TRIAC dimmer or a trailing-edge TRIAC dimmer. In one conventional technology, a rectified output voltage generated by a rectifier is compared with a threshold voltage Vth_on in order to determine a turn-on time period Ton. If the turn-on time period Ton is approximately equal to the duration of a half cycle of the AC input voltage, no TRIAC dimmer is determined to be included in the lighting system; if the turn-on time period Ton is not approximately equal to but is smaller than the duration of a half cycle of the AC input voltage, a TRIAC dimmer is determined to be included in the lighting system. If a TRIAC dimmer is determined to be included in the lighting system, a turn-on voltage slope Von_slope is compared with the threshold voltage slope Vth_slope. If the turn-on voltage slope Von_slope is larger than the threshold voltage slope Vth_slope, the TRIAC dimmer is determined to be a leading-edge TRIAC dimmer; if the turn-on voltage slope Von_slope is smaller than the threshold voltage slope Vth_slope, the TRIAC dimmer is determined to be a trailing-edge TRIAC dimmer.\nIf a conventional lighting system includes a TRIAC dimmer and light emitting diodes (LEDs), the light emitting diodes may flicker if the current that flows through the TRIAC dimmer falls below a holding current that is, for example, required by the TRIAC dimmer. As an example, if the current that flows through the TRIAC dimmer falls below the holding current, the TRIAC dimmer may turn on and off repeatedly, thus causing the LEDs to flicker. As another example, the various TRIAC dimmers made by different manufacturers have different holding currents ranging from 5 mA to 50 mA.\nThe light emitting diodes (LEDs) are gradually replacing incandescent lamps and becoming major lighting sources. The LEDs can provide high energy efficiency and long lifetime. The dimming control of LEDs, however, faces significant challenges because of insufficient dimmer compatibility. For certain historical reasons, the TRIAC dimmers often are designed primarily suitable for incandescent lamps, which usually include resistive loads with low lighting efficiency. Such low lighting efficiency of the resistive loads often helps to satisfy the holding-current requirements of TRIAC dimmers. Hence the TRIAC dimmers may work well with the incandescent lamps. In contrast, for highly efficient LEDs, the holding-current requirements of TRIAC dimmers usually are difficult to meet. The LEDs often need less amount of input power than the incandescent lamps for the same level of illumination.\nIn order to meet the holding-current requirements of the TRIAC dimmers, some conventional techniques use a bleeder for a lighting system. FIG. 2 is a simplified diagram of a conventional lighting system that includes a bleeder. As shown, the conventional lighting system 200 includes a TRIAC dimmer 210, a rectifier 220, a bleeder 224, a diode 226, capacitors 230, 232, 234, 236 and 238, a pulse-width-modulation (PWM) controller 240, a winding 260, a transistor 262, resistors 270, 272, 274, 276, 278 and 279, and one or more LEDs 250. The PWM controller 240 includes controller terminals 242, 244, 246, 248, 252, 254, 256 and 258. For example, the PWM controller 240 is a chip, and each of the controller terminals 242, 244, 246, 248, 252, 254, 256 and 258 is a pin. In yet another example, the winding 260 includes winding terminals 263 and 265.\nThe TRIAC dimmer 210 receives an AC input voltage 214 (e.g., VAC) and generates a voltage 212. The voltage 212 is received by the rectifier 220 (e.g., a full wave rectifying bridge), which then generates a rectified output voltage 222. The rectified output voltage 222 is larger than or equal to zero. The resistor 279 includes resistor terminals 235 and 239, and the capacitor 236 includes capacitor terminals 281 and 283. The resistor terminal 235 receives the rectified output voltage 222. The resistor terminal 239 is connected to the capacitor terminal 281, the controller terminal 252, and a gate terminal of the transistor 262. The gate terminal of the transistor 262 receives a gate voltage 237 from the resistor terminal 239, the capacitor terminal 281, and the controller terminal 252. The capacitor terminal 283 receives a ground voltage.\nAs shown in FIG. 2, the rectified output voltage 222 is used to charge the capacitor 236 through the resistor 279 to raise the gate voltage 237. In response, if the result of the gate voltage 237 minus a source voltage at a source terminal of the transistor 262 reaches or exceeds a transistor threshold voltage, the transistor 262 is turned on. When the transistor 262 is turned on, through the transistor 262 and the controller terminal 254, a current flows into the PWM controller 240 and uses an internal path to charge the capacitor 232. In response, the capacitor 232 generates a capacitor voltage 233, which is received by the controller terminal 244. If the capacitor voltage 233 reaches or exceeds an undervoltage-lockout threshold of the PWM controller 240, the PWM controller 240 starts up.\nAfter the PWM controller 240 has started up, a pulse-width-modulation (PWM) signal 255 is generated. The PWM signal 255 has a signal frequency and a duty cycle. The PWM signal 255 is received by the source terminal of the transistor 262 through the terminal 254. The transistor 262 is turned on and off, in order to make an output current 266 constant and provide the output current 266 to the one or more LEDs 250, by working with at least the capacitor 238.\nAs shown in FIG. 2, a drain voltage at a drain terminal of the transistor 262 is received by a voltage divider that includes the resistors 276 and 278. The drain terminal of the transistor 262 is connected to the winding terminal 265 of the winding 260, and the winding terminal 263 of the winding 260 is connected to the capacitor 230 and the resistor 279. In response, the voltage divider generates a voltage 277, which is received by the controller terminal 256. The PWM controller 240 uses the voltage 277 to detect the end of a demagnetization process of the winding 260. The detection of the end of the demagnetization process is used to control an internal error amplifier of the PWM controller 240, and through the controller terminal 246, to control charging and discharging of the capacitor 234.\nAlso, after the PWM controller 240 has started up, the resistor 274 is used to detect a current 261, which flows through the winding 260. The current 261 flows from the winding 260 through the resistor 274, which in response generates a sensing voltage 275. The sensing voltage 275 is received by the PWM controller 240 at the controller terminal 258, and is processed by the PWM controller 240 on a cycle-by-cycle basis. The peak magnitude of the sensing voltage 275 is sampled, and the sampled signal is sent to an input terminal of the internal error amplifier of the PWM controller 240. The other input terminal of the internal error amplifier receives a reference voltage Vref.\nAs shown in FIG. 2, the rectified output voltage 222 is received by a voltage divider that includes the resistors 270 and 272. In response, the voltage divider generates a voltage 271, which is received by the controller terminal 242. The PWM controller 240 processes the voltage 271 and determines phase angle of the voltage 271. Based on the detected range of phase angle of the voltage 271, the PWM controller 240 adjusts the reference voltage Vref, which is received by the internal error amplifier.\nThe bleeder 224 is used to ensure that, when the TRIAC dimmer 210 is fired on, an input current 264 that flows through the TRIAC dimmer 210 is larger than a holding current required by the TRIAC dimmer 210, in order to avoid misfire of the TRIAC dimmer 210 and also avoid flickering of the one or more LEDs 250. For example, the bleeder 224 includes a resistor, which receives the rectified output voltage 222 at one resistor terminal of the resistor and receives the ground voltage at the other resistor terminal of the resistor. The resistor of the bleeder 224 allows a bleeder current 268 to flow through as at least part of the input current 264. In another example, if the holding current required by the TRIAC dimmer 210 is small and if the average current that flows through the transistor 262 can satisfy the holding current requirement of the TRIAC dimmer 210, the bleeder 224 is not activated or is simply removed.\nAs shown in FIG. 2, the lighting system 200 includes, for example, a quasi-resonant system with a buck-boost topology. The output current 266 of the quasi-resonant system is received by the one or more LEDs 250 and is determined as follows:\n                              I          o                =                              1            2                    ×                                    V              ref                                      R              cs                                                          (                  Equation          ⁢                                          ⁢          1                )            where I0 represents the output current 266 of the quasi-resonant system of the lighting system 200. Additionally, Vref represents the reference voltage received by the internal error amplifier of the PWM controller 240. Moreover, Rcs represents the resistance of the resistor 274.\nFIG. 3 is a simplified diagram showing certain conventional components of the lighting system 200 as shown in FIG. 2. The pulse-width-modulation (PWM) controller 240 includes a dimming control component 300 and a transistor 350. The dimming control component 300 includes a phase detector 310, a reference voltage generator 320, a pulse-width-modulation (PWM) signal generator 330, and a driver 340.\nFIG. 4 shows certain conventional timing diagrams for the lighting system 200 as shown in FIGS. 2 and 3. The waveform 471 represents the voltage 271 as a function of time, the waveform 412 represents the phase signal 312 as a function of time, the waveform 475 represents the sensing voltage 275 as a function of time, and the waveform 464 represents cycle-by-cycle average of the input current 264 as a function of time.\nAs shown by FIGS. 3 and 4, the lighting system 200 uses a closed loop to perform dimming control. The phase detector 310 receives the voltage 271 through the terminal 242, detects phase angle of the voltage 271, and generates a phase signal 312 that indicates the detected range of phase angle of the voltage 271. As shown by the waveform 471, the voltage 271 becomes larger than a dim-on threshold voltage (e.g., Vth_dimon) at time ta and becomes smaller than a dim-off threshold voltage (e.g., Vth_dimon) at time tb. The dim-on threshold voltage (e.g., Vth_dimon) is equal to or different from the dim-off threshold voltage (e.g., Vth_dimoff). The time duration from time ta to time tb is represented by TR, during which the phase signal 312 is at the logic high level, as shown by the waveform 412. The time duration TR represents the detected range of phase angle of the voltage 271.\nDuring the time duration TR, the sensing voltage 275 ramps up and down. For example, during the time duration TR, within a switching period (e.g., TSW), the sensing voltage 275 ramps up, ramps down, and then remains constant (e.g., remains equal to zero) until the end of the switching period (e.g., until the end of TSW).\nThe phase signal 312 is received by the reference voltage generator 320, which uses the detected range of phase angle of the voltage 271 to generate the reference voltage 322 (e.g., Vref). As shown in FIG. 3, the reference voltage 322 (e.g., Vref) is received by the PWM signal generator 330. For example, the PWM signal generator 330 includes the internal error amplifier of the PWM controller 240. In another example, the PWM signal generator 330 also receives the sensing voltage 275 and generates a pulse-width-modulation (PWM) signal 332. The PWM signal 332 is received by the driver 340, which in response generates a drive signal 342 and outputs the drive signal 342 to the transistor 350. The transistor 350 includes a gate terminal, a drain terminal, and a source terminal. The gate terminal of the transistor 350 receives the drive signal 342. The drain terminal of the transistor 350 is coupled to the controller terminal 254, and the source terminal of the transistor 350 is coupled to the controller terminal 258.\nAs shown by the waveform 475, the reference voltage 322 (e.g., Vref) is used by the PWM signal generator 330 to generate the PWM signal 332, which is then used to control the peak magnitude (e.g., CS_peak) of the sensing voltage 275 for each PWM cycle during the time duration TR. For example, each PWM cycle corresponds to a time duration that is equal to the switching period (e.g., TSW) in magnitude. In another example, if the detected range of phase angle of the voltage 271 (e.g., corresponding to TR) becomes larger, the reference voltage 322 (e.g., Vref) also becomes larger. In yet another example, if the detected range of phase angle of the voltage 271 (e.g., corresponding to TR) becomes smaller, the reference voltage 322 (e.g., Vref) also becomes smaller.\nAccording to Equation 1, if the reference voltage 322 (e.g., Vref) becomes larger, the output current 266 (e.g., Io) of the quasi-resonant system of the lighting system 200 also becomes larger; if the reference voltage 322 (e.g., Vref) becomes smaller, the output current 266 (e.g., Io) of the quasi-resonant system of the lighting system 200 also becomes smaller.\nAs shown by FIG. 2, the cycle-by-cycle average of the input current 264 is approximately equal to the sum of cycle-by-cycle average of the output current 266 (e.g., Io) and the bleeder current 268. During the time duration TR, within each switching cycle of the PWM signal 332, the output current 266 changes with time, so the average of the output current 266 within each switching cycle is used to determine the cycle-by-cycle average (e.g., I_PWM_av) of the output current 266 as a function of time. When the time duration TR becomes smaller, the reference voltage 322 (e.g., Vref) also becomes smaller and the one or more LEDs 250 are expected to become dimmer. When the time duration TR becomes too small, the reference voltage 322 (e.g., Vref) also becomes too small and the cycle-by-cycle average (e.g., I_PWM_av) of the output current 266 during the time duration TR becomes smaller than the holding current (e.g., I_holding) required by the TRIAC dimmer 210. In order to avoid misfire of the TRIAC dimmer 210 and also avoid flickering of the one or more LEDs 250, the bleeder current 268 (e.g., I_bleed) is provided in order to increase the cycle-by-cycle average of the input current 264 during the time duration TR. As shown by the waveform 464, the cycle-by-cycle average of the input current 264 during the time duration TR becomes larger than the holding current required by the TRIAC dimmer 210.\nAs shown in FIG. 3, the driver 340 outputs the drive signal 342 to the transistor 350. The transistor 350 is turned on if the drive signal 342 is at a logic high level, and the transistor 350 is turned off if the drive signal 342 is at a logic low level. When the transistor 262 and the transistor 350 are turned on, the current 261 flows through the winding 260, the transistor 262, the controller terminal 254, the transistor 350, the controller terminal 258, and the resistor 274. If the transistor 350 becomes turned off when the transistor 262 is still turned on, the transistor 262 then also becomes turned off and the winding 260 starts to discharge. If the transistor 350 becomes turned on when the transistor 262 is still turned off, the transistor 262 then also becomes turned on and the winding 260 starts to charge.\nAs shown in FIGS. 2-4, the lighting system 200 uses a closed loop to perform dimming control. For example, the lighting system 200 detects the range of phase angle of the voltage 271, and based on the detected range of phase angle, adjusts the reference voltage Vref that is received by the internal error amplifier of the PWM controller 240. In another example, the lighting system 200 provides energy to the one or more LEDs 250 throughout the entire time period of each switching cycle during the time duration TR, which corresponds to the unclipped part of the waveform of the AC input voltage 214 (e.g., VAC).\nAs discussed above, a bleeder (e.g., the bleeder 224) can help a lighting system (e.g., the lighting system 200) to meet the holding-current requirement of a TRIAC dimmer (e.g., the TRIAC dimmer 210) in order to avoid misfire of the TRIAC dimmer (e.g., the TRIAC dimmer 210) and avoid flickering of one or more LEDs (e.g., the one or more LEDs 250). But the bleeder (e.g., the bleeder 224) usually increases heat generation and reduces energy efficiency of the lighting system (e.g., the lighting system 200). Such reduction in energy efficiency usually becomes more severe if a bleeder current (e.g., the bleeder current 268) becomes larger. This reduced energy efficiency often prevents the lighting system (e.g., the lighting system 200) from taking full advantage of high energy efficiency and long lifetime of the one or more LEDs (e.g., the one or more LEDs 250).\nHence it is highly desirable to improve the techniques of dimming control."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus in which a bias of a focus servo circuit is adjusted so that a signal surface of an optical disc may be placed within a focus depth of a laser beam.\n2. Description of the Prior Art\nConventionally, after irradiating a laser beam onto a signal surface of an optical disc, an optical disc player converts light reflected therefrom into an electric signal, and thereby reads out data which is recorded on the optical disc.\nIn the optical disc player known heretofore, it is customary to execute the following servo control.\nFor example, a light beam emitted from a light source is irradiated to the surface of an optical disc, and a return light beam from the optical disc is received by a photo detector.\nA focus error signal and a tracking error signal are detected on the basis of the outputs of focus sensors which are divisions of the photo detector.\nThe focus error signal and the tracking error signal thus detected are supplied to a focus servo circuit and a tracking servo circuit, respectively and thereby execute focus servo control and tracking servo control.\nFIG. 1 is a block diagram for showing a circuit configuration of a conventional optical disc player. As shown in FIG. 1, an optical disc player 1 comprises an optical pickup 2, a tracking servo circuit 5 and a focus servo circuit 6 which receive, via amplifier 3 and 4 respectively, a tracking error signal and a focus error signal obtained on the basis of the difference between the light quantities of individual light receiving elements of a photo detector in optical pickup 2.\nOptical disc player 1 further comprises a tracking driver 7 for driving an actuator of optical pickup 2 under control on the basis of a servo signal from tracking servo circuit 5 to thereby move an objective lens in the tracking direction, a focus driver 8 for driving the actuator of optical pickup 2 under control on the basis of a servo signal from focus servo circuit 6 to thereby move the objective lens in the focusing direction, and an adding circuit 9A for applying a focus bias, which is obtained from the focus bias generator circuit 9, to the focus error signal outputted from amplifier 4.\nOptical pickup 2 has a known structure wherein an objective lens (not shown) is held to be movable biaxially, so that when a tracking coil and a focus coil provided in the actuator are fed with current, the objective lens can be driven biaxially under control in both of the tracking and focusing directions.\nTracking servo circuit 6 is applied with the tracking error signal obtained from optical pickup 2 and amplified by amplifier 3, and then outputs a tracking control signal to driver 7 in accordance with the tracking error signal so as to minimize the tracking error.\nFocus servo circuit 6 is supplied with the focus error signal obtained from optical pickup 2 and amplified by amplifier 4, and then outputs a focus control signal to driver 8 in accordance with the focus error signal so as to minimize the focus error.\nDriver 7 serves to drive the actuator of optical pickup 2 in response to the tracking control signal received from tracking servo circuit 5 and thereby moves the objective lens in the tracking direction to minimize the tracking error.\nMeanwhile, driver 8 serves to drive the actuator of optical pickup 2 in response to the focus control signal received from focus servo circuit 6 and thereby moves the objective lens in the focusing direction to minimize the focus error.\nSince the minimum point of the focus error signal may sometimes fail to coincide with the least jitter point of a reproduced signal, a focus bias obtained from the focus bias generator circuit 9 is applied to the focus error signal for causing the minimum point of the focus error signal to coincide with the least jitter point of a reproduced signal.\nIn optical disc player 1 of the structure mentioned above, the focus error bias is adjusted in the following manner.\nAt the time of assembling optical disc player 1, focusing is performed in an on-state of the focus servo, and an adjusting rheostat 9B incorporated in the focus bias generator circuit 9 is manually operated while observing the RF signal from optical pickup 2 and monitoring the value of the jitter, whereby an optimal focus bias is determined with respect to the individual optical disc player 1.\nHowever, in the optical disc player of the above structure, a time required for determining the focus bias is long with another disadvantage of necessitating the adjusting rheostat 9B, hence raising a problem of higher cost with regard to the component parts.\nFurther, in an operation of recording data on and/or reproducing the same from the optical disc, the refractive index of the optical disc to an incident light beam is rendered different if the material of the disc is different, so that the incident light quantity of the return light beam to the photo detector may also be changed. In addition, occurrence of an ambient temperature fluctuation brings about some harmful influences inclusive of a positional deviation of the objective lens due to the resultant temperature fluctuation in the apparatus.\nAlso, after the focus bias is determined by an experiment and the like, a determined focus bias is stored in a memory and the like which is included in the optical disc player. However, when the bias is thus pre-set, the focus bias is changed due to the aging of the optical disc player and thereby may incur the malfunction of the optical disc player.\nConsequently, it become difficult to adjust the focus bias exactly to its optimal value in the reproduction of data from the optical disc."}
{"text": "Antianxiety drugs, also known as anxiolytic drugs or tranquilizers, are increasingly important psychotropic drugs. They include the benzodiazepines, such as diazepam, chlordiazepoxide, gidazepam, oxazepam, phenazepam, lorazepam, and the like. These benzodiazepines have been the most widely used drugs for the treatment of anxiety. Other antianxiety drugs include, buspirone, gepirone, ipsapirone, and other drugs with high affinity and selectivity for 5HT1.sub.a receptor sites. See for example, Abou-Gharbia et al., J. Med. Chem., 32, 1024 (1989).\nThese agents have been very effective but there has been increasing concern about the disadvantages associated with benzodiazepine therapy. The spectrum of their pharmacological activity, in addition to anxiolytic activity, includes sedative, anticonvulsant, miorelaxant, and amnestic effects, which are often considered both unnecessary and undesirable in the treatment of pathological anxiety. Thus a substantial need exists for nontoxic, highly active, selectively anxiolytic drugs without sedative, muscle relaxant and amnestic activities, which can be used for the treatment of anxiety. The present invention is directed to addressing this need.\nA number of compounds having 1,2,3,4-tetrahydropyrrolo-[1,2-a]-pyrazine heterocycle have been reported. For example, nonsubstituted 1,2,3,4-tetrahydropyrrolo-[1,2-a]-pyrazine is described by A. Skoldinov et al. (U.S. Pat. No. 4,230,856 issued Oct. 28, 1980; USSR Patent No. 798,104 published Jan. 25, 1981). This compound is a precursor in the synthesis of octahydropyrrolo-[1,2-a]-pyrazine, which is useful for the synthesis of physiologically active drugs. 1- and 1,2-substituted 1,2,3,4-tetrahydropyrrolo-[1,2-a]-pyrazine derivatives, having antidepressant activity, are described by I. Jirkovsky (U.S. Pat. No. 4,188,389 issued Feb. 12, 1980). 1-substituted 1,2,3,4-tetrahydropyrrolo-[1,2-a]-pyrazines, having hypotensive activity, are described by V. Peresada et al., Khim-Farm. Zh., 21(9), 1054 (1987). 1-substituted 1,2,3,4-tetrahydropyrrolo-[1,2-a]-pyrazines, possessing coronary-dilating activity, are described by V. Peresada et al., Khim-Farm. Zh., 22(10), 1193 (1988)."}
{"text": "The present invention relates to a biaxially stretch-oriented multilayer film essentially comprising a propylene homopolymer, comprised of a highly pigmented base layer and a thin modified top layer applied to one or both of its surfaces. The invention also relates to the manufacture and use of the film.\nDecorative laminate panels are widely used for producing furniture, e.g., kitchen and office furniture, but also for outdoors applications, e.g., as windowsills, etc. Special types are used under extremely severe conditions, e.g., in shower cubicles or bathrooms or as benches in chemical laboratories. The last mentioned applications, in particular, require panels which are extraordinarily resistant to scratches, to organic and aqueous solvents and to detergents. Apart from imparting a special optical appearance, the surface texture is of particular importance in respect to these characteristic features.\nLaminate boards of the type described above have, for example, been disclosed by German Offenlegungsschrift No. 34 18 282. As a rule, the core layer of a laminate panel is comprised of a multilayer structure comprising a plurality of paper webs impregnated with a phenolic resin. A cellulose paper which is impregnated with a melamine resin and which may be unicolored or printed with a multicolored pattern, is usually applied to this core layer as a top layer.\nPrior to the actual pressing of the top layer(s) and the core layer, the resin-impregnated papers of the core and top layers are dried and thereby partially condensed. The subsequent actual pressing is in general performed by means of heatable or coolable, automatically controlled multi-level pressing units where up to 40 laminate panels can be produced simultaneously. So-called high pressure laminates (H.P.L.) are compressed at pressures of between 70 and 100 bar and temperatures of between 140.degree. and 160.degree. C. During the compression, the resins flow and are cured, whereby a compact, non-meltable, rigid, crosslinked product is formed. Suitable presses include intermittently and continuously working units and through-feed units.\nIf a plurality of panels are stacked on top of one another in the presses, which is economically advantageous for low-thickness core layers, the individual panels must be separated by separators. To produce smooth surfaces, polished metal sheets having plane high-gloss surfaces are used as separators. If the laminate panels are to exhibit a textured or matte surface (for example, resembling veining or fabric texture), appropriately textured templates are required.\nIn order to prevent a sticking-together of the templates and resinous surfaces in those cases where textured templates are used, an additional release sheet must be inserted between the actual separator and the surface of the laminate. Apart from paper and aluminum foil, biaxially stretch-oriented polypropylene foil has been increasingly used for this purpose in recent times.\nWhen clear (transparent) stretch-oriented polypropylene films are employed as release sheets, high-gloss surface finishes are obtained. If, however, matte surface finishes are desired, the manufacturers have to fall back on brushed aluminum foils. More recently, highly pigmented, biaxially stretch-oriented polyester films have also been employed. It is a common disadvantage of these two release films, which are used to produce matte or structured surfaces, that they are relatively expensive.\nAttempts have been made to employ a highly pigmented, i.e., opaque, polyolefin film as a release film, instead of the expensive, highly pigmented polyester films. Polyolefin films of this type are, for example, described in German Offenlegungsschrift No. 28 14 311 (equivalent to U.S. Pat. No. 4,303,708). There is, however, the risk that the surfaces of the laminate panels stick to these release films."}
{"text": "The present invention relates to a flame retardant resin material and a flame retardant resin composition, and more particularly to a flame retardant resin material and a flame retardant resin composition which are improved in flame retardancy, thermal stability or thermal decomposition resistance, and moisture resistance.\nIn order to prevent flame, it is required that resin compositions has a flame retardancy. Usually, halogen flame retardants are used as flame retardants whilst antimony trioxide is used as a co-flame retardant co-used along with the flame retardant. The halogen flame retardants generate harmful halogen substances, typically dioxins. The antimony trioxide as the co-flame retardant has a chronic toxicity. For those reasons, the above substances raise a problem in safety in fire or waste disposal. Phosphoric flame retardants such as red phosphorus and ester phosphate are effective to avoid the above problem. Those phosphoric flame retardants provide influences to moisture resistance of the resin compositions. Particularly, insulators for electronic components are required to have a high reliability. Those phosphoric flame retardants are a problem in use for the insulators for electronic components.\nOn the other hand, epoxy resin compositions are superior in mechanical properties, adhesive property, chemical resistance property, heat resistance and insulating properties, for which reason the epoxy resin compositions are used in various fields in adhesive, coating materials, laminated plates, molding materials and injection materials. In case of the epoxy resin compositions, halogen flame retardants are used as flame retardants whilst antimony trioxide is used as a co-flame retardant. If the flame retardant and the co-flame retardant are used for the epoxy resin composition, problems in not only safety but also corrosion of metals are raised. If those epoxy resin compositions are used as insulators for the electronic components, corrosion resistance to interconnections under high temperature is lowered, whereby reliability of the electronic device is deteriorated. For this reason, it had been required to develop other epoxy resin compositions free from the halogen flame retardant and antimony trioxide.\nIt was investigated to improve the flame retardancy of the resin material by introducing a triazine ring into a molecular structure of an epoxy resin or a phenol resin. In Japanese laid-open patent publication No. 8-311142, it is disclosed that mixtures of phenols with compounds having triazine rings and with aldehydes or phenol condensates such as phenol triazine resins are used as a hardening agent for the epoxy resin compositions. In Japanese laid-open patent publication No. 10-279657, it is disclosed that a phenol triazine epoxy resin obtained by glycidyl-etherification of the above described phenol triazine resin is used as a main component of the epoxy resin composition.\nThere is, however, the following problem in introducing the triazine rings into the molecular structures of the epoxy resins and the phenol resins.\nThe flame retardancy of the resin compositions including the phenol triazine resins and the phenol triazine epoxy resins is exhibited due to a flame reducing mechanism by flame resistant gases which contain, as a main component, nitrogen compounds generated by decomposition of triazines. If in order to emphasize the flame reducing effect, a content of nitrogen in the resin composition is increased, then the resistance to the thermal decomposition of the resin composition is deteriorated, whereby the flame retardancy is thus deteriorated. Since triazines have hydrophilicity, the increase in content of the triazines (nitrogen) in the resin composition causes a remarkable reduction in moisture resistance.\nConsequently, it is difficult to further improve the flame retardancy of the resin composition by introducing nitrogen compounds into the molecular structure of the resin composition.\nIn the above circumstances, it had been required to develop a novel flame retardant resin material and a novel flame retardant resin composition free from the above problems.\nAccordingly, it is an object of the present invention to provide a novel flame retardant resin material free from the above problems.\nIt is a further object of the present invention to provide a novel flame retardant resin material having a high frame retardancy.\nIt is a still further object of the present invention to provide a novel flame retardant resin material having a high thermal stability or a high thermal decomposition resistance.\nIt is yet a further object of the present invention to provide a novel flame retardant resin material having a high moisture resistance.\nIt is further more of the present invention to provide a novel flame retardant resin composition free from the above problems.\nIt is moreover object of the present invention to provide a novel flame retardant resin composition having a high frame retardancy.\nIt is another object of the present invention to provide a novel flame retardant resin composition having a high thermal stability or a high thermal decomposition resistance.\nIt is still another object of the present invention to provide a novel flame retardant resin composition having a high moisture resistance.\nThe present invention provides a flame retardant phenol resin material which includes a phenol condensate, wherein a poly-aromatic compound obtained by a reaction of phenols (A) to aromatics (B) except for phenols and a heterocyclic compound (C) including nitrogen as heteroatom are condensed via aldehydes (D).\nThe present invention also provides a flame retardant epoxy resin material which includes an epoxy resin obtained by glycidyl-etherification of at least a part of phenolic hydroxyl groups of a poly-aromatic compound obtained by a reaction of phenols (A) to aromatics (B) except for phenols and a heterocyclic compound (C) including nitrogen as heteroatom are condensed via aldehydes (D).\nThe present invention also provides a flame retardant resin composition including a flame retardant phenol resin material which includes a phenol condensate, wherein a poly-aromatic compound obtained by a reaction of phenols (A) to aromatics (B) except for phenols and a heterocyclic compound (C) including nitrogen as heteroatom are condensed via aldehydes (D).\nThe present invention also provides a flame retardant resin composition including a flame retardant epoxy resin material which includes an epoxy resin obtained by glycidyl-etherification of at least a part of phenolic hydroxyl groups of a poly-aromatic compound obtained by a reaction of phenols (A) to aromatics (B) except for phenols and a heterocyclic compound (C) including nitrogen as heteroatom are condensed via aldehydes (D).\nThe above and other objects, features and advantages of the present invention will be apparent from the following descriptions."}
{"text": "Field of the Invention\nEmbodiments of the present invention generally relate to performing capacitance sensing while updating a display, or more specifically, to performing capacitance sensing when display updating is paused.\nDescription of Related Art\nInput devices including proximity sensor devices (also commonly called touchpads or touch sensor devices) are widely used in a variety of electronic systems. A proximity sensor device typically includes a sensing region, often demarked by a surface, in which the proximity sensor device determines the presence, location and/or motion of one or more input objects. Proximity sensor devices may be used to provide interfaces for the electronic system. For example, proximity sensor devices are often used as input devices for larger computing systems (such as opaque touchpads integrated in, or peripheral to, notebook or desktop computers). Proximity sensor devices are also often used in smaller computing systems (such as touch screens integrated in cellular phones)."}
{"text": "1. Field of the Invention\nThe present invention relates to a power saving method in a wireless LAN system, and more particularly, to a power saving method in a wireless LAN system conforming to IEEE802.11.\n2. Description of the Related Art\nPower saving in a wireless LAN system conforming to IEEE802.11 is performed in sequences illustrated in FIGS. 1A, 1B, for example, as disclosed in IEEE Std 802.11, 9.7 Frame exchange sequences Table 21 Frame sequence.\nFIG. 1A is a sequence chart illustrating a sequence of operations when no data destined to a terminal station is stored in a base station. The terminal station transmits PS-Poll, which is a control packet for prompting the base station to transmit data, to the base station for requesting the base station to transmit downlink data destined thereto after the terminal station has transitioned to an active mode, in which a transmission/reception function turns on (step 101). Upon receipt of PS-Poll without error, the base station transmits a successful reception notification signal ACK (step 102). The terminal station transitions to a doze state in which the transmission/reception function turns off after it has completed the transmission of ACK.\nFIG. 1B is a sequence chart illustrating a sequence of operations when data destined to the terminal station has been stored in the base station. The terminal station transmits PS-Poll to the base station, after it has transitioned to the active mode, for requesting the base station to transmit downlink data destined thereto (step 111). Upon receipt of PS-Poll without errors, the base station transmits a data frame (step 112). Upon receipt of the data frame without errors, the terminal station transmits ACK (step 113). The terminal station transitions to the doze state after it has completed the transmission of ACK.\nFIG. 2 is a flow chart illustrating the operation of the terminal station in the foregoing situation. As the terminal station starts a receiving operation, it transitions to an active state (step 201). Next, the terminal station transmits PS-Poll to the base station for requesting the base station to transmit data destined to the terminal station (step 202). Next, the terminal station determines whether a response to PS-Pall from the base station is ACK or data (step 203). If the terminal station fails to receive either ACK or data, the terminal station again transmits PS-Pall (step 202). Upon receipt of ACK, the terminal station determines that there is no data destined thereto stored in the base station, and transitions to the doze state (step 205). Upon receipt of data, the terminal station transmits ACK to the base station (step 204), and transitions to the doze state (step 205), followed by termination of the receiving operation.\nIn the foregoing sequence of operations, when data destined to the terminal station is stored in the base station, the base station returns a data frame when it receives the PS-Poll signal from the terminal station without errors. However, conventional general wireless LAN base stations are often configured to return ACK to a terminal station in response to PS-Poll from the terminal station irrespective of the presence or absence of data to the terminal station. This is because it takes long time to complete the transmission of the data frame so that the base station informs, as a temporary measure, the terminal station that the PS-Poll signal has been received without errors. IEEE802.11 also approves that ACK is returned.\nFIGS. 3A, 3B illustrate sequence charts when a base station returns ACK to a terminal station in response to PS-Poll from the terminal station irrespective of the presence or absence of data destined to the terminal station.\nFIG. 3A is a sequence chart when no data destined to a terminal station is stored in the base station. The terminal station transmits PS-Poll to the base station after it has transitioned to the active mode for requesting the base station to transmit downlink data destined thereto (step 301). Upon receipt of PS-Poll without errors, the base station transmits ACK (step 302).\nFIG. 3B is a sequence chart when data destined to a terminal station is stored in a base station. The terminal station transmits PS-Poll to the base station after it has transitioned to the active mode for requesting the base station to transmit downlink data destined thereto (step 311). Upon receipt of PS-Poll without errors, the base station transmits ACK to the terminal station (step 312). When data destined to the terminal station is stored in the base station, the base station further transmits a data frame to the terminal station (step 313). Upon receipt of the data frame without errors, the terminal station transmits ACK (step 314).\nIn this way, data communications in a power save mode in a general wireless LAN differs in a frame sequence chart between the terminal station and base station depending on the presence or absence of data destined to the terminal station stored in the base station. Then, as illustrated in FIGS. 3A, 3B, when the base station returns ACK to the terminal in response to PS-Poll from the terminal station irrespective of the presence or absence of data destined to the terminal station, the terminal station experiences difficulties in determining whether to transition to the doze state when no data is stored in the base station, resulting in difficulties in a reduction in power consumption. This is because the terminal station cannot determine the presence or absence of data destined thereto at the time it receives ACK, and therefore cannot promptly transition to the doze state. This is because the terminal station does not know whether or not data destined thereto is present in the base station at the time it receives ACK. Therefore, even if there is no data destined to the terminal station, the terminal station cannot transition to the doze state immediately after the receipt of ACK."}
{"text": "In recent years, a network facsimile device for communicating facsimile image information by use of electronic mail on the Internet has been developed.\nThe document RFC(Request For Comments) 2301˜2306 published by the organization responsible for collecting technologies relating to the Internet referred to as IETF(Internet Engineering Task Force), prescribes the technical contents of a communication protocol, etc. employed in such a network facsimile device.\nIn the Internet, a delivery system has been developed for delivering electronic mail to a sending terminal for confirming the delivery of electronic mail to a destination terminal. This type of e-mail is hereinafter referred to as “a delivery confirming mail”. This type of delivery confirming mail system is capable of confirming whether electronic mail has been delivered to the address of the communication partner to which it has been addressed. This is an expanded function of the electronic mail system.\nIn such a delivery confirming mail system, the electronic delivery confirming mail is created both when the electronic mail has been delivered to the mail address to which it was sent and when the electronic mail has not been successfully delivered to the mail address to which it has been sent. In each case, electronic delivery confirming mail is created and delivered to the communication partner which sent the electronic mail.\nIn the above-mentioned network facsimile device, however, there presently exists no method of surely knowing whether electronic mail which is used for carrying the image information has been delivered to the address of the communication partner to which it is to be delivered. Accordingly, other ways of confirming the delivery of the image information must be considered.\nHowever, confirming the delivery of image information by a network facsimile device which uses e-mail delivery presents its own unique problems. For example, when the network facsimile device receives electronic mail addressed to itself, the network facsimile device performs local processing on the electronic mail. Consequently, if a delivery confirming mail were used to confirm delivery of the electronic mail sent from the network facsimile device, the network facsimile device would record and output the contents of the delivery confirming mail. However if, all of the contents of the delivery confirming mails are recorded and outputted, there arises the problem of wasteful consumption of recording paper. In addition, a user would need to be present at the network facsimile device to read the delivery confirming mail to determine whether the delivery was successful. These and other problems need to be addressed and solved."}
{"text": "1. The Field of the Invention\nThe present invention generally relates to Reliable Messaging protocols for Web Services. More specifically, the present invention provides for a mechanism that leverages characteristics of Reliable Messaging protocols for Web Services (RM-WS) for verifying and maintaining connection liveliness in a sequence session.\n2. Background and Related Art\nComputer networks have enhanced our ability to communicate and access information by allowing one computer or device to communicate over a network with another computing system using electronic messages. When transferring an electronic message between computing systems, the electronic message will often pass through a protocol stack that performs operations on the data within the electronic message (e.g., parsing, routing, flow control, etc.). The Open System Interconnect (OSI) model is an example of a network framework for implementing a protocol stack.\nThe OSI model breaks down the operations for transferring an electronic message into seven distinct layers, each designated to perform certain operations in the data transfer process. While protocol stacks can potentially implement each of the layers, many protocol stacks implement only selective layers for use in transferring data across a network. When data is transmitted from a computing system, it originates at the application layer and is passed down to intermediate lower layers and then onto a network. When data is received from a network it enters the physical layer and is passed up to the higher intermediate layers and then is eventually received at that application layer. The application layer—the upper most layer—is responsible for supporting application and end-user processing. Further, within the application layer there may reside several other layers (e.g., the Simple Open Access Protocol (SOAP) layer). Another layer incorporated by most protocol stacks is the transport layer. An example of a transport layer is the Transmission Control Protocol (TCP).\nWeb Services (WS) have been a driving force in advancing communications between computing systems and are turning the way we build and use software inside-out. Web Services let applications share data and—more powerfully—invoke capabilities from other applications without regard to how these applications where built; what operating systems or platform they run on; and what devices are used to access them. Web Services are invoked over the Internet by means of industry-standard protocols including SOAP, XML (extensible Markup Language), UDDI (Universal, Description, Discovery and Integration), WSDL (Web Service Description Language), etc. Although Web Services remain independent of each other, they can loosely link themselves into a collaborating group that performs a particular task.\nCurrent WS technologies offer direct SOAP-message communication between an initiator (e.g., a client) and an acceptor (e.g., a service). In the common bi-directional messaging case, a SOAP request message is sent from the initiator to the acceptor and a SOAP reply message is sent in response thereto. Another communication variant between endpoints is unidirectional message exchange, where the initiator sends a message to the acceptor with no response.\nA key benefit of the emerging WS architecture is the ability to deliver integrated, interoperable solutions. Because, however, Web Services provide various services from different business, originations, and other service providers via unreliable communication channels such as the Internet, reliability of WS becomes an increasing important factor. Reliability of WS is impacted by several factors including but not limited to, the reliability of the Web Service end points; reliability characteristics of the communication channel over which the Web Services are accessed; performance and fault-tolerance characteristics; and the extent to which Web Services can handle concurrent client access.\nThere have been attempts at accomplishing reliable messaging of Web Services by choosing a reliable transport protocol over which the messages (e.g., SOAP messages) are exchanged between endpoints. For example, a reliable messaging transport such as message-queues can be used to deliver messages reliably between initiators and acceptors. Messing-queuing communication technologies enable applications on different systems to communicate with each other by sending messages to queues and reading messages from queues that are persisted across failures for reliability.\nAlthough queuing systems offer a transport that can be used to carry SOAP messages reliably, there are several drawbacks to such systems. For instance, these systems offer solutions for an asynchronous operation where the requests (and possibly their responses) are transferred and processed with isolation. Accordingly, these systems are typically heavyweight in terms of resources; involving multiple intermediaries with durable transacted message stores and with considerably more complexity in deployment, programming model and management. All of this is unnecessary for reliable direct communication, and detracts from the goal of minimizing latency. Further, the program model does not directly support request-response style programming or sessions. Accordingly, the queued communication model is different from the current “interactive” Web Services model, and does not address critical “connected” scenarios and “interactive” applications. For example, it is not well suited for cases where a response is expected in a timely manner, or for cases where distributed-transaction-context need to be shared between initiator and acceptor.\nThere have also been attempts at defining reliable transfer layers over fundamentally unreliable transport protocols, e.g., reliable HTTP or HTTPR. A common problem, however, that plagues this solution—as well as the queuing solution—is that reliable messaging can be achieved only if the specific reliable transport protocol is used for communication between the initiator and the acceptor. The fundamental nature of Web Services calls for independence from specific vender platform, implementation languages and specific transport protocols. In a generic case, an initiator may not be able to transmit a message directly to an acceptor using a particular protocol (e.g., acceptor does not support the protocol) or the message may need to pass through multiple hops after leaving the sending node prior to arriving at that destination node. Depending on the nature of the connectivity between the two nodes involved in a particular hop, a suitable transport protocol that does not offer reliable messaging characteristics may have to be chosen.\nIntermediaries may also exist at different levels in the protocol stack; and therefore not offer full end-to-end reliability. For example, transport protocols may offer reliability across lower level intermediaries (e.g., IP level intermediaries—e.g., IP routers). The transport protocol may end, however, at a SOAP intermediary or application layer. Accordingly, the transport protocol may not be able to offer reliability across that intermediary, i.e., no end-to-end reliability across the application layer.\nMore recently, various Reliable Messaging protocols for Web Services (hereinafter referred to as “RM-WS protocols”), e.g., WS-ReliableMessaging, offer solutions to the above identified-deficiencies of current reliable messaging systems. These protocols are transport agnostic connected protocols that allow messages to be delivered reliably between end-point applications in presences of software component, system or network failures. Accordingly, RM-WS protocols offer solutions for reliable, end-to-end, session-oriented communication between an initiator and an acceptor.\nThese RM-WS protocols are akin to TCP in that TCP offers reliable, exactly-once, in-order delivery of a stream of bytes from a TCP sender to TCP receiver across Internet Protocol (IP) routers and multiple networks. Reliable Messaging protocols for WS offer the same and more for messages (note: the unit of transfer is a message, not a byte or a collection of bytes as is the case for TCP wherein the size of the bytes is determined by the available payload space in the IP packet) across multiple intermediaries (including SOAP level intermediaries), transports and connections. Although TCP and RM-WS protocols are both “reliable” protocols, because RM-WS resides at the application or SOAP layer in the OSI model, RM-WS protocols provide for reliable messaging regardless of the transport protocol used to transfer the data. Accordingly, RM-WS protocols are not tied to a particular transport or other protocol used to transfer message between endpoints.\nAlthough a few RM-WS protocols have been around for some time there are still several drawbacks and deficiencies of these protocol specs. For example, these RM-WS protocols do not provide a way for verifying the connectivity of an established sequence session. Further, the specifications do not provide for a way to maintain the liveliness of a connection and/or to extend the sequence session in the absence of an exchange of messages over an inactivity timeout period. Accordingly, there exists a need to verify and maintain connection liveliness for systems that use a RM-WS protocol."}
{"text": "In parametric internal model-based control (IMC), a controller may be used to control a process where the controller includes a model of the process. The accuracy of the model drives the performance of the controller. IMC is also known in the art as model predictive control (MPC) or simply as model-based control (MBC). For purposes of this disclosure, the term IMC will be used.\nFIG. 1 shows a representation of a control system 10 that includes a model 12 of an integrating process 14 that may be controlled with a controller 16 by a predictive control technique that is represented by an inverted version of the model 12 (realizable model inverse) and an associated filter. Parameters of the model are used to calculate, or predict, a future process variable (PV), and generate a control output (CO) responsive to a given set point. But, as will be explained in greater detail below, modeling the behavior of an integrating process that may have one or more poles at origin leads to computational difficulties.\nFor any non-zero disturbance (e.g., input disturbance d1 and/or output disturbance d2), the predicted process variable (also known as a controlled variable or CV) will grow without bound due to the integrating nature of the conventional implementation of the model 12 and associated inverted model. The same holds for any non-zero control output CO (also known as a manipulated variable) as the model prediction of the future state of PV also will grow without a bound.\nThe model output under the forgoing conditions is represented in the graphs of FIGS. 2a and 2b, which share the same scale for the time axis. In FIG. 2a, the control output CO is graphed. As illustrated, a step change to the control output CO is made at a given time. As will be appreciated, the control output CO change is the input value to the process model 12 and the process 14. The predicted process variable that is output by the model 12 is graphed in FIG. 2b. \nA dead time D is a measurement of how long the process takes to start to respond to the CO change and may be defined by the amount of time that elapses from the change in control point to when the predicted process variable exceeds a noise threshold defined by a noise band. A time constant (or lag time) may be used as an indicator of how long the process takes to reach steady state after the dead time. For an integrating process, the time constant value may be the amount of time that elapses from the end of the dead time to when the slope of the predicted process variable achieves 63% of the slope of the steady state PV. 95% of the slope of the steady state PV may be achieved after three time constants.\nAs illustrated, using an integrating process model, the value for the predicted process variable will grow without bound. A maximum slope of the predicted process variable curve (or PV maximum slope) may be determined and the gain (k) may be defined by the PV maximum slope divided by the CO change. The unbounded increase in predicted process variable is not necessarily a design issue for the control system 10, but represents a computational and numerical issue.\nThree common solutions exist, but each have drawbacks. The first solution is to convert output disturbances (or modeling errors) to an equivalent input disturbance, such that disturbances and errors that are input to the model are mathematically zeroed in the steady state. The second solution is to reinitialize the state of the predicted process variable from time to time to reduce the occurrence of a numerical overflow. The third solution is to use an internally stable algorithmic internal model control (AIMC) of a two degree of freedom IMC.\nWhile these solutions improve the behavior of the model, they all introduce design and computational complexities. Moreover, for an integrating process, the design of an internal model-based controller is complicated by the introduction of a higher order filter as will be demonstrated below.\nWith continued reference to FIG. 1, conventional representations of a first order integrating process will now be described. Under these representations, the integrating process 14 may be represented by equation 1A and the model 12 of the process may be represented by equation 2.\n                                          G            p                    ⁡                      (            s            )                          =                              k            s                    ⁢                      ⅇ                          -              DS                                                          Eq. 1A                                                      G            m                    ⁡                      (            s            )                          =                                            k              ^                        s                    ⁢                      ⅇ                                          -                                  D                  ^                                            ⁢              S                                                          Eq. 2            \nThe steady state process gain is represented by k in equation 1A and an estimated steady state process gain is represented by {circumflex over (k)} in equation 2. Similarly, D is the process dead time and {circumflex over (D)} is an estimated dead time. In each equation, a Laplace variable is represented by S.\nThe controller 16 may be represented by equation 3, in which Gm+−1 is an invertible portion of the model 12 (equation 2) and F(s) is a filter that is designed to make the controller represented by Gc(s) realizable. The filter specifies a desired response for the process variable PV (referred to as a desired trajectory of PV).Gc(s)=Gm+−1·F(s)  Eq. 3\nA controller computation round-off error is shown in FIG. 1 as d3, which represents a rounding error. While the rounding error may be very small, the rounding error may still lead to a numerical instability issue, as described below.\nThe system 10 of FIG. 1 may be simplified into the forms illustrated in FIGS. 3 and 4. If, in FIG. 3, a perfect model 12 is obtained for a first order process 14, Gm(s) will equal Gp(s), and then both Gm(s) and Gp(s) may be represented by equation 1A. If a first order filter F(s) equaling 1/(εs+1) is used, equation 3A results, where ε is a filter time constant.\n                              G          c                =                                            G                              m                +                                            -                1                                      ·                          F              ⁡                              (                s                )                                              =                      s                          k              ⁡                              (                                                      ɛ                    ⁢                                                                                  ⁢                    s                                    +                  1                                )                                                                                  Eq          .                                          ⁢          3                ⁢        A            \nAs indicated, Gm+ is the invertible portion of the model and is equal to k/s. In FIG. 4, GIMC may be considered the combination of Gc and Gm and represents an internal model-based controller 16′. GIMC may be expressed as equation 4.\n                              G          IMC                =                                            G              c                                      1              -                                                G                  c                                ⁢                                  G                  m                                                              =                                                    s                                  k                  ⁡                                      (                                                                  ɛ                        ⁢                                                                                                  ⁢                        s                                            +                      1                                        )                                                                              1                -                                                      ⅇ                                          -                      DS                                                                                                  ɛ                      ⁢                                                                                          ⁢                      s                                        +                    1                                                                        =                          s                              k                ⁡                                  (                                                            ɛ                      ⁢                                                                                          ⁢                      s                                        +                    1                    -                                          ⅇ                                              -                        DS                                                                              )                                                                                        Eq        .                                  ⁢        4            \nIn equation 5A, the process variable PV is expressed using a three term equation that represents a dynamic relationship among the value at the input to the controller 16′ (first term where SP is the set point), the value at the input to the process 14 following introduction of an input disturbance d1 (second term) and the value at the output of the process following introduction of an output disturbance(s) d2 (third term).\n                              PV          ⁡                      (            s            )                          =                                                                              G                  IMC                                ⁢                                  G                  p                                                            1                +                                                      G                    IMC                                    ⁢                                      G                    p                                                                        ⁢            SP                    +                                                    G                p                                            1                +                                                      G                    IMC                                    ⁢                                      G                    p                                                                        ⁢                          d              1                                +                                    1                              1                +                                                      G                    IMC                                    ⁢                                      G                    p                                                                        ⁢                          d              2                                                                    Eq          .                                          ⁢          5                ⁢        A            \nThe second term of equation 5A corresponds to an input disturbance to the process 14 and it is desirable for this term to be zero at steady state. If the second term is not zero, a steady state error occurs. By applying a step input for the set point (or set point change) an analysis of how PV will behave may be made. One may represent the step input as r, such that SP(s) equals (1/s)r and equation 6 follows from equation 5A.\n                                                                        t                ⁢                                  →                  lim                                ⁢                                  ∞                  ⁢                                                                          ⁢                                      PV                    ⁡                                          (                      t                      )                                                                                  =                              s                ⁢                                  →                  lim                                ⁢                                  0                  ⁢                                                                          ⁢                                                                                    G                        IMC                                            ⁢                                              G                        p                                                                                    1                      +                                                                        G                          IMC                                                ⁢                                                  G                          p                                                                                                      ⁢                  r                                                                                                        =                              s                ⁢                                  →                  lim                                ⁢                                  0                  ⁢                                                                          ⁢                                                                                    ⅇ                                                  -                          DS                                                                                            (                                                                              ɛ                            ⁢                                                                                                                  ⁢                            s                                                    +                          1                          -                                                      ⅇ                                                          -                              DS                                                                                                      )                                                                                    1                      +                                                                        ⅇ                                                      -                            DS                                                                                                    (                                                                                    ɛ                              ⁢                                                                                                                          ⁢                              s                                                        +                            1                            -                                                          ⅇ                                                              -                                DS                                                                                                              )                                                                                                      ⁢                  r                                                                                                        =                                                s                  ⁢                                      →                    lim                                    ⁢                                      0                    ⁢                                                                                  ⁢                                                                  ⅇ                                                  -                          DS                                                                                                                      ɛ                          ⁢                                                                                                          ⁢                          s                                                +                        1                        -                                                  ⅇ                                                      -                            DS                                                                          +                                                  ⅇ                                                      -                            DS                                                                                                                ⁢                    r                                                  =                r                                                                        Eq        .                                  ⁢        6            \nIt may be observed that for a change in the set point SP variable, there is a direct correlation in the process variable PV attained at steady state. But if a step input disturbance (e.g., d1(s) equaling d1/s) is introduced, a steady state error in PV results as indicated by equation 7.\n                                                                        t                ⁢                                  →                  lim                                ⁢                                  ∞                  ⁢                                                                          ⁢                                      PV                    ⁡                                          (                      t                      )                                                                                  =                              s                ⁢                                  →                  lim                                ⁢                                  0                  ⁢                                                                          ⁢                                                            G                      p                                                              1                      +                                                                        G                          IMC                                                ⁢                                                  G                          p                                                                                                      ⁢                                      d                    1                                                                                                                          =                              s                ⁢                                  →                  lim                                ⁢                                  0                  ⁢                                                                          ⁢                                                                                    k                        ⁢                                                                                                  ⁢                                                  ⅇ                                                      -                            DS                                                                                              s                                                              1                      +                                                                        ⅇ                                                      -                            DS                                                                                                    (                                                                                    ɛ                              ⁢                                                                                                                          ⁢                              s                                                        +                            1                            -                                                          ⅇ                                                              -                                DS                                                                                                              )                                                                                                      ⁢                                      d                    1                                                                                                                          =                              s                ⁢                                  →                  lim                                ⁢                                  0                  ⁢                                                                          ⁢                                                            k                      ⁡                                              (                                                                              ɛ                            ⁢                                                                                                                  ⁢                            s                                                    +                          1                          -                                                      ⅇ                                                          -                              DS                                                                                                      )                                                                                    s                      ⁡                                              (                                                                              ɛ                            ⁢                                                                                                                  ⁢                            s                                                    +                          1                          -                                                      ⅇ                                                          -                              DS                                                                                +                                                      ⅇ                                                          -                              DS                                                                                                      )                                                                              ⁢                                      d                    1                                                                                                                          =                                                s                  ⁢                                      →                    lim                                    ⁢                                      0                    ⁢                                                                                  ⁢                                                                                            k                          ⁢                                                                                                          ⁢                          ɛ                                                +                                                  kD                          ⁢                                                                                                          ⁢                                                      ⅇ                                                          -                              DS                                                                                                                                                                            ɛ                          ⁢                                                                                                          ⁢                          s                                                +                        1                        +                                                  ɛ                          ⁢                                                                                                          ⁢                          s                                                                                      ⁢                                          d                      1                                                                      =                                                                            (                                                                        k                          ⁢                                                                                                          ⁢                          ɛ                                                +                                                  k                          ⁢                                                                                                          ⁢                          D                                                                    )                                        ⁢                                          d                      1                                                        ≠                  0                                                                                        Eq        .                                  ⁢        7            \nIt may further be observed that for a step output disturbance where d2(s) equals d2/s, the third term of equation 5A advantageously goes to zero as demonstrated by equation 8.\n                                                                        t                ⁢                                  →                  lim                                ⁢                                  ∞                  ⁢                                                                          ⁢                                      PV                    ⁡                                          (                      t                      )                                                                                  =                              s                ⁢                                  →                  lim                                ⁢                                  0                  ⁢                                                                          ⁢                                      1                                          1                      +                                                                        G                          IMC                                                ⁢                                                  G                          p                                                                                                      ⁢                                      d                    2                                                                                                                          =                                                s                  ⁢                                      →                    lim                                    ⁢                                      0                    ⁢                                                                                  ⁢                                                                                            ɛ                          ⁢                                                                                                          ⁢                          s                                                +                        1                        -                                                  ⅇ                                                      -                            DS                                                                                                                                                ɛ                          ⁢                                                                                                          ⁢                          s                                                +                        1                        -                                                  ⅇ                                                      -                            DS                                                                          +                                                  ⅇ                                                      -                            DS                                                                                                                                              =                0                                                                        Eq        .                                  ⁢        8            \nIt may be concluded that with a first order filter and a perfect model, the IMC controller 16′ (FIG. 4) will have a steady state error for a step input disturbance. To reduce or eliminate the steady state error for a step input disturbance, a complex filter F(s) may be employed. But such a filter is difficult to implement in an actual controller used to control an actual integrating process. In addition, the model that drives the controller still contains a numerical issue in that the value output by the model will grow without bound due to the exponential component of the integrating model."}
{"text": "An extremely fine pattern of a solid element such as large-scale semiconductor integrated circuit or the like is formed by using chiefly a reduction projection exposure method that is one of optical lithography methods. This method is a method of reducing and transferring a mask pattern formed by a photomask or a reticle (called a mask hereinafter), onto a substrate by use of an imaging optical system.\nImprovement of resolution in the reduction projection exposure method is advanced by high numerical aperture in the imaging optical system and short wavelength of exposed light. However, since there are needs of more extreme fineness for least process size of the solid element than the above-mentioned improvement, a deformed illumination exposure method or a phase shift mask exposure method, so-called a super resolution exposure method is developed and applied.\nThe phase shift mask exposure method includes, for example, a Levenson type phase shift mask, a halftone type phase shift mask, an auxiliary pattern arrangement type phase shift mask, and the like. The Levenson type phase shift mask is a mask generating a phase difference of 180 degrees between light beams that permeate regions between adjacent apertures (light permeating region) on the mask. The Levenson type phase shift mask also has an effect on further improvement of the resolution thereof within regions in which pattern-arranged pitches are extremely fine. For example, if a Levenson type phase shift mask is used, a reduction projection exposure method using KrF excimer laser light can eminently improve resolving characteristics even within a size region less than a least process size whose sufficient resolution is difficult to obtain in the case of use of normal masks. Further, the halftone type phase shift mask is a mask in which a halftone film is formed on a mask substrate instead of a light shield film. The halftone film has functions of making exposed light beams be permeated some per cents and of generating a phase difference of 180 degrees between exposed light beams permeating the halftone film and permeating apertures around which the halftone film is removed.\nAnd, the auxiliary pattern arrangement type phase shift mask is a mask having such a size as not to resolve on a semiconductor wafer around a main aperture and arranging auxiliary patterns for generating a phase difference of 180 degrees between exposed light beams permeating the main aperture. The auxiliary pattern arrangement type phase shift mask can be used when mask patterns are not arranged densely. For example, in a mask pattern for transferring isolated hole patterns, there is a structure of arranging auxiliary patterns which have such a size as not to be transferred on the semiconductor wafer of a plane surface containing upper, lower, right and left side of the main aperture and which generate a shift difference of 180 degrees relative to exposed light beams permeating the main apertures. This results in improvement of a light intensity profile of the main aperture and enhancement the resolving characteristics. This method is described in Japanese Patent Laid-open No. 5-19446, which discloses a technique of disposing auxiliary patterns on an end of dense patterns and around isolated patterns in order to enhance resolution of the dense patterns end and the like. Further, for example, Japanese Patent Laid-open No. 6-123963 discloses a technique of disposing respective auxiliary patterns such that light beams permeating respective adjacent patterns do not interfere with one another, or a technique of disposing one auxiliary pattern relative to the main aperture when the auxiliary patterns are arranged between the adjacent patterns. And, for example, Japanese Patent Laid-open No. 6-289591 discloses a technique of disposing auxiliary patterns in a symmetrically shifted manner in order to enhance flexibility in arrangement of the main apertures. Further, for example, Japanese Patent Laid-open No. 8-297359 discloses a technique of making layouts of mask patterns such that one main aperture and one auxiliary pattern are handled as one unit in order to facilitate the layouts of the mask patterns. And, for example, Japanese Patent Laid-open No. 11-84625 discloses a structure of disposing main apertures, auxiliary patterns, and shifters arranged like zigzag at dense main apertures, and of arranging the auxiliary patterns at each end of memory mats."}
{"text": "Known optical instruments for measuring angles, such as sextants and octants, comprise support means in the form of a frame; an eyepiece mounted on the support means; a first reflector, known as a \"horizon glass\", mounted on the support means in alignment with the axis of the eyepiece; an adjustment member, known as an \"index bar\", pivotally mounted on the support means; a movable reflector mounted on the adjustment member so as to face towards the first reflector; and measuring means comprising an angularly graduated \"arc\" for measuring movement of the adjustment member. The support means comprise first and second arms extending divergently from a hub and interconnected at their free ends by an arcuate limb. The eyepiece and first reflector are respectively mounted on the first and second arms and the adjustment member carrying the movable reflector is pivotally mounted in the hub at the junction of the first and second arms.\nTo use an instrument of this type to measure the inclination of a celestial body, such as the sun or a star, it is necessary to hold the instrument with the optical axis of the eyepiece aligned with the horizon so that the horizon may be viewed through the eyepiece and a non-reflective portion of the horizon glass and then, with the pivotal axis of the hub arranged horizontally, above the optical axis of the eyepiece, to measure the angular movement of the adjustment member from a datum position in which the movable reflector is parallel to the first reflector to a position in which light from the celestial body is reflected from the movable reflector to the first reflector and back along the optical axis of the eyepiece so that the celestial body can be viewed through the eyepiece. This angular movement is equal to half the angle of inclination of the celestial body from the axis of the eyepiece. This is consistent with the law of optics whereby the angle between the first and last directions of a beam of light which is reflected twice, by two plain reflective surfaces, is twice the angle between the two reflective surfaces. By reference to navigational tables, it is possible to establish from such inclinations the positions from which measurements have been made.\nAlthough an octant, in which the first and second arms are inclined at 45.degree., is capable of measuring 45.degree. movement of the adjustment member and therefore 90.degree. of inclination from an axis extending from the eyepiece to the horizon, this is not sufficient for measuring all inclinations if the instrument is used by an operator standing on an elevated platform such as the bridge of a ship. In this case, the axis extending from the eyepiece to the horizon dips below horizontal and so the maximum elevation which can be measured is less than vertical. To overcome this disadvantage, it is necessary to use the more expensive sextant in which the first and second arms are inclined at 60.degree. and which is capable of measuring inclinations of up to 120.degree. from the optical axis of the eyepiece.\nIn any case, to be of value, conventional sextants and octants have to be accurately made and are thus often too expensive for all but the most serious and professional navigators. Moreover, these instruments are particularly difficult to use, and cannot be used at all when the horizon is obscured by haze or at night. The elevations of stars must therefore be measured during twilight, at dusk and dawn, when both the stars and the horizon are visible. Clearly, on very many occasions when the horizon is indistinct, the great precision of navigational instruments such as sextants and octants is unnecessary in view of the unavoidable error in aligning the optical axis of the eyepiece with the horizon.\nOne known technique for modifying a conventional sextant or octant so as to make it possible to measure inclinations when the horizon is not visible is to provide the instrument with a spirit level and complex optical means which enable an operator to view an artificial horizon when the instrument is held so that the optical axis of the eyepiece extends horizontally. However, the resultant instrument, known as a \"bubble sextant\", does not work satisfactorily and has never been widely accepted by navigators. The modification increases the complexity of the instrument with attendant difficulty in keeping the instrument in accurate operating adjustment, and the cost is considerably higher than for conventional sextants and octants."}
{"text": "The present invention relates to an X-Y input device, and more particularly to an X-Y input device suitable for use as an input device associated with a graphic display apparatus.\nGraphic display apparatus basically include a display screen, a display controller, a data channel, and an input device which may be of various types. One known input device is a \"joystick\" having a lever supported by a gimbal mechanism and tiltable by the operator in any direction. A control device detects the direction and angle of tilt of the lever and generates voltages or digital signals indicative of coordinate values in X and Y directions. This type of input device is disadvantageous however in that the range of angular movement of the lever is limited and data signals entered by the operator are relatively unstable.\nIn an effort to eliminate the above shortcomings, there has in recent years been developed an input device called a \"mouse.\" One type of the mouse has a rotatable member such as a steel ball, and first and second driven rollers held in contact with the ball are rotated in response to rotation thereof. The first and second driven rollers have their axes of rotation extending substantially perpendicularly to each other. The mouse also includes first and second angle detector means often comprised of variable resistors or encoders for separately detecting the angles of rotation of the first and second driven rollers. The ball, first and second driven rollers, and first and second angle detector means are all housed in a casing.\nThe casing has an opening defined in its bottom with the ball partly projecting through the opening. In use, the casing is held by the operator to place the ball against a given base or surface. By moving the casing to cause the ball to roll on the surface in any desired direction, the first and second driven rollers are rotated about their own axes through angles dependent on the rolling movement of the ball. The directions and angles of rotation of the driven rollers are converted by the first and second angle detector means into voltages or digital signals corresponding to the rolling movement of the ball representative of coordinate values in X and Y directions. The generated signals are then entered into a display apparatus."}
{"text": "Increasing amounts of available petroleum feedstocks correspond to natural gas sources or other methane-containing sources. The increased availability of methane and/or small hydrocarbon petroleum sources can potentially have increased value if efficient methods can be identified for conversion to larger compounds.\nOne option for converting methane (and other small hydrocarbons) to other compounds is to reform the methane to form H2 and/or CO (i.e., synthesis gas). Both steam reforming and dry reforming are known, but each type of reforming poses a variety of challenges. In particular, both types of reforming processes are prone to formation of coke on the reforming catalyst. This is especially true when the natural gas feed contains elevated levels of CO2 and/or C2+ hydrocarbon molecules, such as ethane or propane. In such cases, conventional reforming processes employ additional gas separations processes to reduce the CO2 content to acceptable levels—about 18-20% for steam reforming processes and much lower, about 5%, for autothermal reforming processes—or a prereformer to convert C2+ hydrocarbons to carbon monoxide and hydrogen.\nThese pretreatments increases the capital and operating expenses of reformning processes. Moreover, natural gas streams with high amounts of CO2 and/or C2+ hydrocarbons are more likely to be used as low cost fuels rather than for the production of high margin products through their conversion into synthesis gas, and are thus cost suppressed.\nU.S. Pat. No. 8,454,911 describes methods of using a reverse flow reactor for reforming of methane to form acetylene.\nU.S. Pat. No. 7,217,303 describes methods for pressure swing reforming of a hydrocarbon fuel in the presence of steam to form hydrogen. The reforming step is described as having a peak temperature of 700° C. to 2000° C.\nU.S. Pat. No. 7,815,873 describes methods for reforming of a hydrocarbon fuel in a reverse flow reactor configuration. The reforming step is described as having a temperature of at least 1000° C."}
{"text": "This disclosure relates to computer animation and computer generated imagery. More specifically, this disclosure relates to systems and methods for creating and using hierarchies of models for rigging.\nComputer-generated imagery (“CGI”) is a process of using computers to convert specifications of objects, things, lights, effects, etc. into images. CGI might be used for generating one image, for example, of a teapot with one light source, or might be used to generate an animation, a sequence of images such as might be used for a feature-length film.\nMany computer graphic images are created by mathematically modeling the interaction of light with a three dimensional scene from a given viewpoint. This process, called rendering, generates a two-dimensional image of the scene from the given viewpoint, and is analogous to taking a photograph of a real-world scene. Animated sequences can be created by rendering a sequence of images of a scene as the scene is gradually changed over time. A great deal of effort has been devoted to making realistic looking rendered images and animations.\nSome animation might be done by manually painting each image, but with the wide use of computers, it is common, for movies or other features, to have a user (e.g., an animator or other skilled artist) specify geometric descriptions of models or other objects, such as characters, props, background, or the like that may be rendered into images. An animator may also specify poses and motions for objects or portions of the objects. In some instances, the geometric description of objects may include a number of animation variables (avars), and values for the avars. As an example, a simple character might comprise a model that comprises a fixed body, two arms and two legs, with four joints where the arms and legs join the body and can be moved. An animation variable for such a model might be the angle between the body and the right arm. The animator could “pose” this model by placing the body at a location and an orientation in the virtual three-dimensional space, then specifying the four angles (two arm rotations and two leg rotations). This posing might be done by the animator typing in values for each degree of freedom and viewing a rendering of the result.\nIn such a simple case, entering in a few values is not a problem, but many models are much more complicated. For example, a character model might include a movable body, arms, elbows, fingers, facial muscles, etc., all with their own degrees of freedom. To simplify the entry of the animator's desired pose and/or movements, a user interface with “handles” or “widgets” for avars might be provided. As an example, a user interface might display a rendered character and overlay yellow squares that the animator can drag and move to indicate movements desired by the animator. Internally, the animation editing software that presents the user interface to the user and receives the animator's inputs could translate a “drag-handle” operation into a change in value of an avar.\nIn some animation terminology, a model refers to a collection of elements that together form an object that is presented in a scene and may have some degrees of freedom. The number of degrees of freedom may depend on the type of object. A model for a building might have a degree of freedom as to how a door opens or closes and where the building is located, but that model—except perhaps in some stories—might not be expected to have a degree of freedom that allows a central spine to rotate and bend.\nA rigging for a model refers to the user interface (logical or actual) for exercising the degrees of freedom. Using the rigging, an animator can specify a pose of the model for a still image and/or specify how the model is to move from pose to pose in an animation. Rigging is not part of the displayed scene, although there are user interfaces provided to animators that illustrate parts of the rigging, such as when animation is being edited. Models can be given various controllers, animation variables, and handles for an animator to manipulate various locations of the object's topology to create complex movements and motions.\nFor some models, a bone/joint system can be set up to deform various locations of the object's topology. For example, the bone/joint system can be connected to foot, ankle, knee, hip, and other leg locations of a humanoid model to provide the structure to make the humanoid model walk. Other types of information may also be “hung” on the object's topology to add further realism or additional control for the animator. In other words, information may be associated with a vertex, edge, span, or face of the mesh that forms to the object's topology. However, the above processes can be very involved and time consuming to simply generate a single model.\nAdditionally, a typical feature-length animation may require hundreds to thousands of models. This increases the production time and cost of the animation if each model may be required to be hand-created and setup. One possible solution can be to hand copy the information from one model to another. However, this process still requires an animator to place or “hang” the copied data onto the correct position of the new objects topology. Rarely is each character exactly the same, so each character's topology can have some differences that the animator has to deal with.\nFor example, an animated feature film might involve classes of objects or characters, with hierarchical subclasses. For example, there might be class of characters such as humans, with subclasses for warriors, elderly humans, tall humans, short stocky humans, children, adults, etc. and subclasses of subclasses, such as injured tall humans, etc. It could be that each of the different subclasses have separate models and separate rigging, but that can be a lot of work for an animation team.\nAccordingly, what is desired are improved methods and apparatus for solving some of the problems discussed above, while reducing further drawbacks, some of which are discussed above."}
{"text": "Permanent seat licenses, also known, for instance, as personal seat licenses, Stadium Builder licenses (SBL), Charter Ownership agreements (COA), Charter Seat licenses (CSL), and Preferred Access Speedway Seating (PASS), are a relatively recent innovation wherein, for example, sports teams sell “seat licenses” to fans desiring admission tickets to team events. Generally, the seat license for a particular seat within the stadium or arena enables and/or requires the holder of that license to purchase season tickets to team events. The season ticket, in turn, gains the holder admission to individual team events during the course of the season where the holder is able to sit in the particular seat. However, permanent seat licenses (referred to herein as “PSL's”) are also an important source of revenue for the sports team issuing the PSL's. This revenue can be used by the team to, for example, subsidize the construction of new facilities or for other purposes requiring capital outlay. Note that the PSL's may be issued by entities other than just sports teams. For example, the sports team may often act in partnership with a government entity, where the partnership issues the PSL's. In other instances, the government entity may build, own, and operate the stadium independently of the sports team and may thus issue the PSL's independently thereof. Further, the venue may be owned by an entity separate from both the sports team and the government entity and it is the separate entity which issues the PSL's. Note further that the PSL concept as described herein may be applicable to venues for other events besides sports such as, for example, theaters or the like. Thus, it will be understood by one skilled in the art that the PSL concept may have wide and varied applicability.\nPSL's are generally offered directly to the public by the sports team (referred to herein as the “issuer”). These PSL's often have prices ranging from hundreds to many thousands of dollars. Sometimes, for a variety of reasons, such as moving from the area or losing interest in the team, holders may sometimes wish to sell their PSL's after buying them from the issuer. A problem herein lies that, since a PSL is a niche commodity, there is no established resale market for the product. In some instances, computer bulletin boards and newspaper classified ads are utilized to find a subsequent buyer. However, such transactions are often laden with uncertainty for both the seller and the buyer. For instance, due to the cost of the PSL's, a buyer may often have to procure financing, which may delay the transaction. If the seller is desperate to sell the PSL's, he may arrange to finance all or part of the sale to the buyer, which may further lead to uncertainty for the seller recouping his money from the buyer. Likewise, the buyer is faced with uncertainty as to, for example, whether the seller is the rightful owner of the PSL's or whether the seller will deliver the PSL's upon the buyer paying the purchase price. Sometimes, the buyer may be faced with the uncertainty that the PSL's correspond to seats that do not have the view of the playing field that the buyer desires or that the seats corresponding to the PSL's have a view that may adversely affect the value of the PSL's. In other instances, the PSL's may be offered for sale by the seller independently of the current season's tickets, wherein the buyer would not be able to obtain tickets until the following season. This situation may often not be revealed to the buyer, or the buyer may forget to inquire about this detail, until the transaction has been completed.\nThe PSL transaction may further include mutual uncertainties for both the buyer and the seller. For instance, since there is no established market, it is often difficult for both the seller and the buyer to ascertain a going market price for the PSL's. Such a market price may also be affected by performance factors with the team itself. For example, PSL's may rise significantly in price if the team has come off a successful championship season or has recently built a new stadium or arena. Conversely, the value of a PSL may drop significantly if the team is on a losing streak or is otherwise forecast to perform poorly in the future. In addition, there may be instances where the transaction between the buyer and the seller occurs over a distance and the parties cannot meet in person. In these instances, the uncertainties facing both the buyer and seller as to the reliability of the transaction may take on a greater significance. Further, the issuer typically maintains records of the current ownership of the respective PSL and the seller and/or the buyer may be required to notify the issuer of the ownership change at the time of the transaction, as well as to pay an associated fee for the transfer. Issuer requirements with regard to the transaction may vary between issuers and may also comprise a procedure which is unknown to the parties.\nThus, there exists a need for a medium for buying and selling PSL's in which reasonable assurances are provided to both the buyer and the seller concerning important aspects of the transaction, for example, that the seller is the rightful owner of the PSL and that the buyer is capable of paying the purchase price. Such a medium should also desirably be able to inform the parties of at least a best estimate of a current market price for the PSL's, be able to let the buyer evaluate the views from the corresponding seats, and be able inform the buyer and seller as to the transaction requirements imposed by the respective issuer. It may also be desirable for the medium to incorporate provisions for facilitating the transaction wherein the seller can submit the PSL and the buyer can submit the corresponding remittance to a trusted intermediary. The medium should also desirably be able to account for the current season's tickets corresponding to the PSL together with or independently of the sale of the PSL."}
{"text": "A common cause of failure of implanted biomedical devices is infection. The attachment of bacteria to medical implants and in-dwelling catheters, and the proliferation of such bacteria, is a major cause of infection during or after the implantation process. Treating an implant with antibiotics has not proven very effective to combat infections, and has sometimes resulted in the development of resistant strains of bacteria.\nSilver coatings are effective to combat infections; however, silver coatings tend to cause tissue irritation, largely due to excessively rapid release of the silver into the surrounding tissues.\nAn antimicrobial coating for a medical implant is needed which has hardness and low friction properties, and that will result in a slow, long-term release of antimicrobial metal atoms into the body in order to prevent infection."}
{"text": "The flexible Ethernet (FlexE) combines some technical features of the Ethernet and a transport network (for example, an optical transport network (OTN), and a synchronous digital hierarchy (SDH)), and is an important milestone in the evolution of an Ethernet technology. With emergence of a flexible Ethernet technology, Ethernet physical interfaces show virtualization characteristics. A plurality of Ethernet physical interfaces may be bonded together, to support several virtual logical ports. For example, a 400-gigabit (400G) flexible Ethernet physical interface group obtained by bonding four 100-gigabit Ethernet (100GE) physical interfaces may support several logical ports.\nThe Ethernet physical interface is an asynchronous communications interface, and is allowed to have a clock frequency difference of ±100 ppm (one ten-thousandth). For example, in 10GE, for two physical interfaces whose nominal bandwidths are 10 G, one bandwidth may be one ten-thousandth larger than the nominal value, and the other bandwidth is one ten-thousandth smaller than the nominal value, that is, 10G*(1+0.0001) and 10G*(1−0.0001). A clock frequency at the logical port inherits a clock frequency characteristic on the physical interface, and therefore the logical port also has a difference of 100 ppm. For example, actual bandwidths of two logical ports that are formed by different physical interfaces or physical interface groups and whose nominal bandwidths are 25 G may be approximately 25G*(20460/20461)*(1+0.0001) and 25G*(20460/20461)*(1−0.0001) when overheads of timeslot division and timeslot management in the flexible Ethernet are considered. When the flexible Ethernet is used to bear a service, idle-code-block (Idle) insertion or deletion needs to be performed hop by hop, to adapt a service rate to a bandwidth rate difference between the service and the physical interfaces or the logical ports. FIG. 1 is a schematic diagram of service transport in the flexible Ethernet in the prior art. As shown in FIG. 1, when a service between customer devices Ca and Cb is borne by using flexible Ethernet devices Pa, Pb, and Pc, the Pa, the Pb, and the Pc needs to perform idle-unit insertion or deletion.\nHowever, idle-code-block insertion or deletion causes loss of a clock frequency and time phase information of the service, that is, the clock frequency and the time phase information of the service cannot be transparently transported, and consequently the clock frequency and the time phase cannot be synchronized between a source network device and a sink network device of the service."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus such as a copying machine, a printer and a facsimile machine, and more particularly, to an image forming apparatus which corrects skew feeding of a sheet while conveying the sheet, and which conducts positioning of a side edge (end) of the sheet in a width direction intersecting with a sheet conveying direction.\n2. Description of Related Art\nGenerally, image forming apparatuses of an electrophotographic system, an offset printing system and an ink-jet system are known. As image forming apparatuses using the electrophotographic system, image forming apparatuses of various systems are known, such as an image forming apparatus of a direct transfer system which transfers a toner image from a photosensitive drum directly to a sheet, and an image forming apparatus of an intermediate transfer system which once transfers a toner image to an intermediate transfer member and then transfers the toner image to a sheet. As image forming apparatuses using the electrophotographic system, an image forming apparatus of a tandem system in which a plurality of image forming portions is arranged, and an image forming apparatus of a rotary system in which a plurality of image forming portions is cylindrically arranged are known.\nIn recent years, in image forming apparatuses of the electrophotographic system, making full use of a merit that a plate is not formed, apparatuses targeting a printing market of small copies are provided. In order to be accepted by such a light printing market, high velocity (high productivity) and high image quality must be achieved in various kinds of materials, and a requirement for sheet conveying precision is increased. An image position precision with respect to a sheet is most required to be high, an image position deviation of front and back when images are formed on both surfaces is also included. There is a method for adjusting a position of an image with respect to a sheet, but a method for adjusting a sheet with respect to an image predominates.\nThe precision of an image position is determined by registration of a sheet in a sheet conveying direction, registration of a sheet in a width direction intersecting with the sheet conveying direction, magnification and skew feeding. Among them, it is difficult to correct the skew feeding of a sheet by electrical control. For example, if skew feeding of a sheet is detected and an image which inclines corresponding to the skew feeding is formed, it is possible to correct an image position with respect to the sheet. However, in the case of a color image on which three or four colors are superposed, if an image is inclined every sheet, color is changed in every sheet due to deviation of dot formation of each color. Further, since it takes time to calculate for inclination of an image, productivity is largely reduced. Thus, the skew feeding of a sheet is determined by performance of conveying precision of a sheet.\nGenerally, the skew feeding of a sheet and registration are independently controlled, but in recent years, there is proposed a method for correcting the skew feeding and correcting registration in a direction intersecting with the sheet conveying direction at the same time or by the same driving operation (see Japanese Patent Laid-Open No. 10-310289). More specifically, it includes two moving drive motors which independently slide two rollers arranged in the sheet conveying direction in a direction intersecting with the sheet conveying direction, and two optical sensors which detect a side edge (end) of a sheet are arranged in the sheet conveying direction corresponding to the roller. Control is carried out such that the rollers are slid in the width direction so that the side edge of the sheet follows the optical sensor.\nAccording to the conventional method, however, since the rollers are slid in the direction intersecting with the sheet conveying direction to correct the skew feeding of the sheet or correct the position of the sheet in the width direction, stress is applied to the sheet when the rollers slide. Especially, in the case of a thin paper sheet, since a sheet is bent between the two rollers, it is difficult to precisely correct the skew feeding of a sheet or correct a position of a sheet in the width direction.\nAccording to the conventional method, since the moving drive motor is rotated normally and reversely by an ON/OFF operation of the optical sensor, a sheet overshoots in the width direction and reciprocates, and it takes time for a side edge of the sheet to reach a target position. If an attempt is made to enhance the position precision of an image with respect to a sheet, the conveying velocity of the sheet cannot be increased, and the productivity cannot be enhanced.\nHence, the present invention provides an image forming apparatus which can handle various kinds of materials including a sheet having weak elasticity such as a thin paper sheet, and which has excellent position precision of an image with respect to a sheet, and which enhances the productivity."}
{"text": "Pipelines, such as oil, water or sewer distribution or collections systems, are constructed by welding together a plurality of pipes, often at the installation site of the pipeline.\nKnown in the art are automated orbital pipeline welding systems, such as disclosed in U.S. Pat. No. 4,373,125, issued Feb. 8, 1983 to Kazlauskas. These automated systems are capable of forming accurate and strong welds on large diameter pipes. Such systems are relatively large, typically weighing over 500 kilograms, and are installed at a stationary location, such as an oil rig. Accordingly, such systems are not suitable for mobile field use.\nMore mobile welding systems are known. Typically, such systems comprise a welding carriage, or \"bug\" as often termed in the art, which includes a welding torch. The torch may be suited for Gas Metal Arc Welding (GMAW), Shielded Metal Arc Welding (SMAW) or Gas Tungsten Arc Welding (GTAW). The carriage is typically the size of a hand-held power tool and is mounted, via constrained rollers, on an annular rail or guide disposed on the weldment. The carriage includes a driving pinion which meshes with a toothed rack of the annular guide, thereby providing a means for guiding the carriage and torch around the weldment. Typically, the carriage has at least one d.c. brush type motor mounted thereon for driving the carriage on the guide as well as a motor for oscillating the torch. Typically too, the electronic control system for controlling the motors is housed in a separate unit, remote of the carriage, and linked thereto through a plurality of control cables.\nThere are various types of control systems for these types of bugs. One type of control system, based on analog electronics, employs potentiometers which are adjustable by an operator during the welding process. These potentiometers control the speed of the carriage drive mechanism motor, the torch oscillating motor, and, if present, a consumable electrode wire feed motor. Operators using such a system almost always adjust the speeds of the various motors during the welding process, particularly in order to avoid the problem of having the deposited weld bead, which is liquid, drip due to the influence of gravity. The problem with using such systems, however, is that welders have complete control of the welding process and can adjust the speeds of the motors such that the resulting weld does not always fall within the requisite specifications for the weld. The problem is further compounded by the fact that often the potentiometers are not linear.\nOther known bug control systems employ largely digital control systems wherein, in combination with suitable carriages, many of the weld parameters can pre-programmed. For example, one known type of mobile automated welding system allows the weld current, arc voltage, welding speed, oscillation speed, width and dwell time, torch height, tilt angle and annular position, to be digitally programmed. This system also provides a programmed means for controlling the carriage travel speed to deal with the deposited metal drip problems. Similar systems known in the art, such as disclosed in U.S. Pat. No. 5,534,676 issued Jul. 9, 1996 to Rinaldi et. al., have more sophisticated methods for accomplish this objective. However, one limitation common to these types of systems is a lack of flexibility in enabling the welder to vary the pre-programmed parameters during the welding process.\nIn any event, these mobile welding machines are often used in some of the harshest and most remote environments in the world. Thus, reliability of the machines is important. There are a number of limitations in the present design of mobile welding machines of the types described above that affect their reliability. Welding machines having the known fully automated digital control systems tend to have many sensors and other delicate mechanisms which are prone to breakage in use, particularly under heavy use in harsh construction environments. Welding machines having the analog control systems require frequent recalibration, particularly under operating conditions wherein the ambient temperature fluctuates widely. In addition, irrespective of the type of control system, brush-type motors mounted on the carriage have a tendency to bum out within a relatively short period of time. Moreover, the signals carried by such cables can be prone to electromagnetic interference caused by nearby operating machinery, particularly high frequency inverter type power sources which radiate relatively large amounts of electromagnetic energy.\nIn addition to having a reliable welding system, it is also important to ensure the quality of the resulting weld, particularly as the weld is being formed. Thus, it is desired to have a real time weld monitoring system. Some of the welding machines of the prior art having automated digital control systems provide a feedback to a remote computer indicating what the actual values of some of the carriage and welding parameters are. However, these systems do not inform the operator in real time whether the weld is being properly made. It would be helpful to have more comprehensive weld quality information readily available so that the operator could immediately adjust certain operating parameters to ensure the quality of the weld.\nThe present invention seeks to address many of the limitations of the prior art mobile pipeline welding systems described above."}
{"text": "Display tables may be used, for example, in a retail setting to present products to consumers and to showcase features of those products."}
{"text": "For various reasons, it is highly desirable to be able to qualitatively and/or quantitatively measure human luteinizing hormone (hLH) and human chorionic gonadotropin (hCG) in a convenient, reliable manner. For example, it is known that the concentration of hLH in female biological fluids increases dramatically just prior to ovulation. A convenient, qualitative detector of the hLH surge for \"at home\" use could pinpoint the time of ovulation for women with fertility problems and/or aid in fertility regulation. Also, a quantitative hLH assay, which could be used by the physician in the office would be useful in the medical evaluation and observation of various conditions associated with hLH presence and/or concentration variation.\nSimilarly, significant presence of hCG in the female biological fluids has been accepted as one of the most reliable tests confirming pregnancy. hCG is secreted by the developing blastocyst and can be detected in pregnant women as early as 7 to 9 days after fertilization. Again, a qualitative \"at home\" pregnancy test capable of providing a reliable result within the first two weeks of pregnancy would be extremely useful. Also, a convenient assay procedure usable by the practicing physician would prove of benefit in the diagnosis and evaluation of various abnormal conditions characterized by hCG levels, such as in the diagnosis of ectopic pregnancy and spontaneous abortion, and in the follow-up of patients with trophoblastic disease and infertility problems.\nOne aspect of the present application is concerned with improved solid phase assay systems and methods for measuring hLH and/or hCG, which do not require radioactive substances. In one embodiment of the present invention, an enzymatic marker is employed. In another embodiment, a direct dye marker is used to eliminate the step of reacting enzyme with its substrate for measurement. Another aspect of the present application is concerned with materials and reagents to be used in the improved assay systems and methods.\nVan Weeman et al, \"Immunoasssay Using Antigen-Enzyme Conjugate\", FEBS LETTERS, 15:232 (1971), conjugated hCG to horseradish peroxidase through glutaraldehyde and then used the conjugate for enzyme-immunoassay of hCG. A solid phase assay procedure is disclosed by Van Weeman et al where hCG antibody is attached to a cellulosic support (reprecipitated and diazotized m-aminobenzyloxymethyl cellulose or microcrystalline cellulose activated by CnBr.) Also, see U.S. Pat. No. 3,654,090 using the Van Weeman conjugate in similar assay procedures. Another enzyme used in the art in place of horseradish peroxoidase is alkaline phosphatase (ALP). For example, see Engvall et al, \"Enzyme-linked Immunoabsorbent Assay (ELISA), Quantitative Assay of Immunoglobulin G\", Immunochemistry, 8:871 (b 1971). Kawaoi et al, \"An Improved Method of Conjugation of Peroxidase with Protein\". Fed. Proc., 32 Abstract 840 (1973) also disclose a method for forming an enzyme-hormone complex.\nSaxena et al, \"Development of a Solid-Phase Centrifugation-free Enzyme Assay for LH for Ovulation Detection\", Psychoneuroendocrinology in Reproduction, Zichella et al, editors, Elsevier/North Holland (1979), p. 277, describe both an enzymeimmunoassay and an enzymereceptorassay for hLH using antibody-coupled glass beads or receptor-coupled glass beads, respectively (glutaraldehyde activated aminopropyl glass beads) and an hLH-alkaline phosphatase conjugate (4-azidobenzoyl derivative of hLH). The receptor used by Saxena et al was of relatively low yield and suffered from poor stability.\nCarlsson et al, \"Protein-Thiolation and Reversible Protein-Protein Conjugation\", Biochem. J., 173:723 (1978), disclose that N-succinimidyl 3-(2-pyridyldithio) propionate (SPDP) can be employed as a heterobifunctional reagent in forming protein-protein conjugates through a disulfide link. As one example, Carlsson et al form a horseradish peroxidase-rabbit anti- (human transferrin) antibody conjugate. In general, 2-pyridyl disulfide structures are introduced into both the peroxidase and the antibodies by their reaction with SPDP. Although either of the 2-pyridyldisulfide moieties can be converted into the corresponding thiol derivative by specific reduction with dithiothreitol at pH 4.5, Carlsson et al formed the thiolated antibodies and reacted them with the peroxidase 2-pyridyl disulfide derivative by thiol disulfide exchange to produce the peroxidase enzyme-antibody conjugate.\nIt is also known that hLH and hCG consist of two noncovalently linked polypeptide chains classified as .alpha. and .beta. subunits. The .alpha. subunits of these glycoprotein hormones, as well as the .alpha. subunits of other glycoprotein hormones, have almost identical amino acid sequences and are practically indistinguishable immunochemically. In contrast, the .beta. subunits of each have their own distinctive amino acid sequences which differentiate them from the .beta. subunits of the other glycoproteins. Methods are known to purify antibody, such as hCG .beta. antiserum, to reduce cross-reactivity with hCG .alpha. subunits or with similar hLH subunits. For example, see Jibiki et al, \"A Receptor-Immunoassay for the Determination of the Specificity of Anti-hCG-.beta. Sera\", ACTA Endocrinologica, 87, 838 (1978), where hCG-.beta. antisera is purified by immunoabsorption with hLH and hCG-.alpha..\nSaito et al, \"Use of Receptors in the Preparation of LH-Free Serum\", J. Clin. Endocrinol. Metab. 43:1186 (1976), used purified gonadotropin receptor to bind serum hLH to the exclusion of other serum proteins. This general technique provides a purer serum, depending upon the specificity of the receptor.\nGlass beads have been used as the solid phase in the radioimmunoassay of hormones such as hLH and hCG. See Post et al, \"A Rapid, Centrifugation-Free Radioimmunoassay Specific for Human Chorionic Gonadotropin Using Glass Beads as Solid Phase\", J. Clin. Endoclinol. Metab., 50:169 (1980), where solid glass beads of 6 mm diameter were sandblasted and then heated in the presence of .gamma.-aminopropyltriethoxysilane to generate reactive alkylamino groups on the glass surface. The amino groups are activated with glutaraldehyde and subsequently hCG-.beta. antibody is covalently coupled thereto.\nFinally, a key reagent in an hLH or hCG receptor assay is the receptor itself. Heretofore, large quantities of the common hLH and hCG receptor of the corpora lutea have not been obtained in sufficiently pure form for use in enzyme receptor assays of good reliability and accuracy. See Khan et al, \"Use of Purified Gonadotropin Receptor in the Development of an Enzyme Receptorassay (ERA) for Luteinizing Hormone (LH) and Human Chorionic Gonadotropin (hCG)\", Enzyme Labelled Immunoassay of Hormones and Drugs, S. B. Pal,editor, de Gruyler & Co. (1978).\nA number of the steps used in the receptor purification system of the present invention, the present invention resulting in a good yield of a receptor having better stability and higher purity, were disclosed by the inventors in the symposium paper. \"Current Status of the Purification and Characterization of LH-hCG Receptors,\" Functional Correlates of Hormone Receptors in Reproduction, Mahesh et al, editors, page 397, Elsevier North Holland (1980), presented Oct. 15, 1980.\nU.S. Pat. No. 4,016,250 and as a divisional thereof, U.S. Pat. No. 4,094,963, both to Saxena, disclose receptorassay methods for determining hCG and/or hLH in a biological fluid wherein binding of said hCG and hLH in the sample is with a plasma membrane extract from the corpus luteum of a species possessing the common receptor for hCG and hLH. The additional purification steps disclosed herein, as compared to the Saxena patents, result in the preparation of a receptor material of significantly higher purity and binding capacity."}
{"text": "As electronic devices become increasingly compact and portable, the need to make precise electrical connections between components on a very small scale increases. Arrays of fine, closely spaced, parallel electrical conductors are often needed to connect arrays of closely spaced, side-by-side pads found on printed circuit boards, liquid crystal displays, display panels, charge-couple devices, or the like.\nIn electronic devices, electronic components are typically arranged with extreme space restrictions. Therefore, electrical connectors are often required to be flexible and to have closely spaced, parallel conductive stripes provided on a flexible insulating support. See, e.g., U.S. Pat. No. 5,059,262 and U.S. Pat. No. 4,931,598. Other methods of producing parallel conductive stripes include: a) conductive inks, b) thin metal wires, or c) stripes of thin metal films, e.g., deposited through a mask or selectively etched to provide the desired conductor width and spacing.\nThere are disadvantages inherent in the presently known electrical connector tape constructions. Electrical connector tapes that use conductive inks as the conductive stripes typically have an undesirably high resistance. The manufacture of electrical connector tape using thin metal wires for the conductive stripes requires drawing the metal wires down to size and attaching the wires to a flexible insulating support. Known methods are typically difficult add expensive. Stripes made by means of photolithographic techniques are often complicated and also expensive.\nIn presently known electrical connector tapes, the tape is often constructed such that the spacing of the individual conductive stripes is the same as that of the terminal pads to which the tape is intended to be bonded. Accordingly, when the bonds are made between the conductive stripes and an array of terminal pads, it is necessary that absolute registration be maintained between the stripes and the pads during bonding. The fine pitch of many arrays of terminal pads makes such registration very difficult. Thus, it is often necessary to use magnifying devices when bonding the electrical connector tape to the terminals. However, if the pitch of the electrical connectors is so fine that one or more conductive stripes will contact a terminal pad during bonding, absolute registration may not be necessary. However, it will still be necessary, in most applications, to maintain a generally parallel alignment to prevent cross-over connections."}
{"text": "The present invention relates to a dust radiation monitor apparatus used in radiation source handling facilities such as nuclear power plants to measure the radioactivity concentration of dust in the air in the facilities, and a dust sampling apparatus used in the dust radiation monitor apparatus.\nConventionally, in radiation source handling facilities such as nuclear power plants, dust radiation monitor apparatuses for measuring the radioactivity concentration of dust in the air in the facilities are widely used.\nFIG. 1 is a schematic view showing an example of the arrangement of a dust radiation monitor apparatus widely used in practice in radiation source handling facilities such as nuclear power plants.\nIn the conventional dust radiation monitor apparatus in FIG. 1, when air is drawn into a chamber C through a sampling pipe 1, radiation dust in the air is collected through a filter (filter paper or the like) in a dust collection section 2. The air after dust collection is exhausted through a sampling pump 3. A .beta. ray detection section 5 detects .beta. rays from the radiation dust collected through the filter, and converts it into an electrical signal. When a .beta. ray measurement value is obtained by a .beta. ray measuring section 7, a data processing section 8 compares the measurement value with a warning set value and determines contamination if any.\nNatural radioactive substances exist in nature. Radon is a naturally occurring substance (to be referred to as a natural nuclide hereinafter). Since radon exists in a rare gas form in nature, it also floats in the air in the facilities. A nuclide of this kind emits .alpha. and .beta. rays in the process of decaying into a daughter nucleus and granddaughter nucleus.\nThe radiation dust collected by the above dust radiation monitor apparatus therefore contains natural nuclides. For this reason, as shown in FIG. 2, a total .beta. ray measurement value Va includes a contribution V2 of natural nuclides as well as a contribution V1 of measurement nuclides due to leakage. Owing to the influences of such natural nuclides, many radioactive nuclides (artificial contamination nuclides) as measurement targets may be determined to exist in amounts larger than actual amounts. For this reason, the influences of natural nuclides must be separately evaluated.\nAs a technique of solving this problem, a technique of accurately determining the presence/absence of radioactive contamination without separately evaluating the influences of natural nuclides is disclosed in Jpn. Pat. Appln. KOKAI Publication No. 9-211133. According to this technique, .alpha. and .beta. rays emitted from natural nuclides are separately measured in advance by a measurement system, and the emission ratio of the measured .alpha. and .beta. rays from natural nuclides is obtained in advance. In actual contamination determination, radiation emitted from measurement targets is detected by a radiation detector, and .alpha. and .beta. rays are separately measured from the detection signal. A correction processing means then obtains a .beta. ray value base on natural nuclides which are contained in the .beta. ray measurement value on the basis of the emission ratio of .alpha. and .beta. rays and the .alpha. ray measurement value, and subtracts the .beta. ray value base on natural nuclides from the .beta. ray measurement value. As a result, a .beta. ray value free from the influences of natural nuclides can be obtained.\nIn the dust radiation monitor apparatus disclosed in Jpn. Pat. Appln. KOKAI Publication No. 9-211133, however, .alpha. and .beta. rays are measured by using a single radiation detector. More specifically, detection/measurement is performed by a single radiation detector having a detection section made up of two layers (.alpha. and .beta. ray detection layers). For this reason, the following problems are posed.\n(1) Low-energy .beta. rays are absorbed by the .alpha. ray detection layer.\nMore specifically, with the structure formed by stacking the .alpha. and .beta. ray detection layers on each other, .beta. rays output from the dust collection section always reach the .beta. ray detection layer through the a ray detection layer. The detection efficiency of low-energy .beta. rays, which are easily absorbed, decreases, resulting in a measurement error.\n(2) Separate measurement of .alpha. and .beta. rays has its own limitation.\nMore specifically, light components from the .alpha. and .beta. ray detection layers mix with each other. The mixed light components are separated by a subsequent circuit in accordance with differences in rise characteristics and emission amount, but they are not perfectly separated. Since .alpha. and .beta. rays as measurement targets vary in energy, they mix with each other in a certain energy region, resulting in a measurement error. Especially, the low-energy side of a rays tend to mix with the high-energy side of .beta. rays.\nDemands therefore have arisen for a dust radiation monitor apparatus which has a structure that prevents light components from .alpha. and .beta. ray detection layers from mixing with each other, and can eliminate any measurement error due to mixing of light components. Demands has also arisen for a dust radiation monitor apparatus which can prevent the other detection layer from absorbing .beta. rays in detecting .beta. rays, and improve the detection efficiency of low-energy .beta. rays, thereby eliminating any measurement error due to .beta. ray absorption.\nIn such a dust radiation monitor apparatus, the following problems arise in intermittent measurement and continuous measurement, in addition to the above problems of a decrease in detection efficiency and measurement errors.\nFIG. 3 is a schematic view showing an example of the arrangement of a conventional intermittent dust radiation monitor apparatus.\nReferring to FIG. 3, a pipe switching unit 1401 sequentially switches connections between the radiation monitor side and a plurality of (n) sampling pipes 1400 which are installed in different sampling places in radiation source handling facilities (not shown) to introduce air from the respective sampling places. As shown in FIG. 3, the pipe switching unit 1401 is constituted by a plurality of solenoid valves 1402.\nThe air introduced through the sampling pipes 1400 sequentially switched by the pipe switching unit 1401 is drawn by a pump 1406 through a pipe system and continuously sent to a dust collection section 1411. Dust in the air sent to the dust collection section 1411 is collected on filter paper 1403 driven by a filter paper driving section 1407.\nIn addition, the amount of air drawn is adjusted to a predetermined amount by a flow rate indicator 1404 and a flow rate control valve 1405.\nRadiation from the dust collected on the filter paper 1403 is detected by a radiation detector 1410, and the radioactivity concentration of the dust on the filter paper 1403 is then measured by a data converter 1408 using the radiation reading from the radiation detector 1410. The measurement result is output to a display/recording section 1409.\nSuch an intermittent dust radiation monitor apparatus can perform only intermittent measurement. For this reason, the flow rate of air must be increased to a value equal to or more than a detection limit value, and the flow rate of air must be fixed to maintain a high detection precision.\nUnder the circumstances, a continuous dust radiation monitor apparatus is proposed, in which dust collection/measurement units are arranged for the respective sampling pipes to perform continuous measurement.\nIn such a continuous dust radiation monitor apparatus, however, since dust collection/measurement units equal in number to the sampling pipes must be installed, the apparatus arrangement becomes large in size.\nDemands have therefore arisen for a dust radiation monitor apparatus which can switch intermittent measurement to continuous measurement to perform continuous monitoring, as needed, while intermittent measurement is performed in normal operation as in the prior art, thereby improving the measurement precision. In addition, demands have arisen for a dust radiation monitor apparatus which can collectively perform dust connection control and measurement processing for each of intermittent and continuous dust radiation monitors by intensively using one data processing section, and can change the schemes of the respective monitors according to the circumstances.\nIn dust sampling apparatus used in the above dust radiation monitor apparatus, the following problems arise in collecting dust with filter paper.\nFIG. 4 shows the arrangement of a conventional dust sampling apparatus.\nA chamber 901 is used to draw air from a predetermined place in a radiation management area and exhaust the air after dust collection. A filter paper holder 905 for holding filter paper 903 and a paper filter receiving wire net 904 is mounted in the chamber 901. The filter paper holder 905 is mounted on the chamber 901 through a paper filter holder O-ring 918. A radiation detector 902 for detecting radiation emitted from the dust collected by the filter paper 903 is mounted in the chamber 901. The radiation detector 902 is held by a detector holder 906, which is mounted on the chamber 901, through a detector O-ring 917.\nIn the dust sampling apparatus having this arrangement, when air is drawn through the inlet of the chamber 901, dust in the air is collected on the filter paper 903 in the filter paper holder 905, and the air after dust collection is exhausted from the outlet of the chamber 901.\nIn this conventional dust sampling apparatus, to increase the detection sensitivity, the distance between the filter paper 903 and the radiation detector 902 is reduced to several millimeters. Owing to this structure, the flow path of air passing through the filter paper 903 is not uniform, and dust is nonuniformly collected on the filter paper 903. As a result, the detection efficiency decreases due to self-absorption of dust, or a sensitivity calibration deviation occurs with respect to a calibration ray source.\nUnder the circumstances, demands have arisen for a dust sampling apparatus which can make the flow path of air passing through filter paper uniform, and can uniformly collect dust, thereby preventing a decrease in detection efficiency due to self-absorption of dust and a sensitivity calibration deviation with respect to a calibration ray source."}
{"text": "1. Field of the Invention\nThis invention relates to synthetic oil-based preferably polyol ester-based turbo oils which use a synergistic combination of phosphorous (P)-based and sulfur (S)-based load additive chemistries which allows the turbo oil formulation to impart high load-carrying capacity and also to meet or exceed US Navy MIL-L-23699 requirements including Oxidation and Corrosion Stability and Si seal compatibility.\nLoad additives protect metal surfaces of gears and bearings against uncontrollable wear and welding as moving parts are heavily loaded or subjected to high temperatures. Incorporating high load-carrying capacity into a premium quality turbo oil without adversely impacting other properties can significantly increase the service life and reliability of the turbine engines.\nThe mechanism by which load additives function entails an initial molecular adsorption on metal surfaces followed by a chemical reaction with the metal to form a sacrificial barrier exhibiting reduced friction between the rubbing metal surfaces. In the viewpoint of this action, the effectiveness as load-carrying agent is determined by the surface activity imparted by a polar functionality of a load additive and its chemical reactivity toward the metal; these features can lead to severe corrosion if not controlled until extreme pressure conditions prevail. As a result, the most effective load additives carry deleterious side effects on other key turbo oil performances: e.g., corrosion, increased deposit forming tendency and elastomer incompatibility."}
{"text": "The semiconductor integrated circuit (IC) industry has experienced rapid growth. In the course of IC evolution, functional density (i.e., the number of interconnected devices per chip area) has generally increased while geometry size (i.e., the smallest component (or line) that can be created using a fabrication process) has decreased. This scaling down process generally provides benefits by increasing production efficiency and lowering associated costs. Such scaling down has also increased the complexity of processing and manufacturing ICs and, for these advances to be realized, similar developments in IC processing and manufacturing are needed."}
{"text": "1. Field of the Invention\nThis invention relates to devices for drilling and boring through subterranean formations. More specifically, this invention relates to polycrystailine diamond compacts (\"PDCs\"), also known as cutting elements or diamond inserts, which are intended to be installed as the cutting element of a drill bit to be used for boring through rock for many applications, including oil, gas, mining, and/or geothermal exploration, that require drilling through geological formations. Still more specifically, this invention relates to polycrystalline diamond inserts which have a surface topography formed integral to an otherwise spherical, conical, or other uniform geometric shape, to increase stress at the insert/rock interface, thereby inducing the rock to fail with the expenditure of less overall energy while introducing little additional internal stresses to the insert.\n2. Description of Related Art\nThree types of drill bits are most commonly used in penetrating geologic formations. These are: (1) percussion bits; (2) rolling cone bits, also referred to as rock bits; and (3) drag bits, or fixed cutter rotary bits. Each of these types of bits may employ the polycrystalline diamond inserts of this invention as the primary cutting device.\nIn addition to the drill bits discussed above, polycrystalline diamond inserts may also be used with other down hole tools, including but not limited to: reamers, stabilizers, and tool joints. Similar devices used in the mining industry may also use this invention.\nPercussion bits penetrate through subterranean geologic formations by an extremely rapid series of impacts. The impacts may be combined with simultaneous rotations of the bit. An exemplary percussion bit is shown in FIG. 1b. The reader is directed to the following list of related art patents for further discussion of percussion bits.\nRolling cone bits currently make up the largest number of bits used in drilling geologic formations. Rolling cone bits have as their primary advantage the ability to penetrate hard geologic formations while still being generally available at a relatively low cost. Typically, rolling cone bits operate by rotating three cones, each oriented substantially transverse to the bits axis and in a triangular arrangement, with the narrow end of each cone facing a point in the direct center of the bit. An exemplary rolling cone bit is shown in FIG. 1a.\nA rolling cone bit cuts through rock by the crushing and scraping action of the abrasive inserts embedded in the surface of the rotating cone. These abrasive inserts are generally composed of cemented tungsten carbide, but may also include polycrystalline diamond coated cemented tungsten carbide insert of this invention, where increased wear performance is required.\nThe primary application of this PDC invention is currently believed to be in connection with percussion and rolling cone bits, although alternative embodiments of this invention may find application in connection with other drilling tools.\nA third type of bit is the drag bit, known also as the fixed cutter bit. An example of a drag bit is shown in FIG. 2. The drag bit is designed to be rotated about its longitudinal axis. Most drag bits employ PDCs which are brazed into the cutting blade of the bit. The PDCs then shear the rock as the bit is rotated about its longitudinal axis.\nIt is expected that this invention will find primary application in percussion and rolling cone bits, although some use in drag bits may also be feasible.\nA polycrystalline diamond compact (\"PDC\"), or cutting element, is typically fabricated by placing a cemented tungsten carbide substrate into a refractory metal container (\"can\") with a layer of diamond crystal powder placed into the can adjacent to one face of the substrate. The can is then covered. A number of such can assemblies are loaded into a high pressure cell made from a soft ductile solid material such as pyrophyllite or talc. The loaded high pressure cell is then placed in an ultra-high pressure press. The entire assembly is compressed under ultra-high pressure and ultra-high temperature conditions. This compression causes the metal binder from the cemented carbide substrate to become liquid and to \"sweep\" from the substrate face through the diamond grains and to act as a reactive liquid phase promoting the sintering of the diamond grains. Sintering of the diamond grains cause the formation of a polycrystalline diamond structure. As a result the diamond grains become mutually bonded together to form a diamond mass over the substrate face. The metal binder may remain in the diamond layer within the pores of the polycrystalline structure or, alternatively, it may be removed via acid leeching and optionally replaced by another material forming so-called thermally stable diamond (\"TSD\"). Variations of this general process exist and are described in the related art. This detail is provided so the reader may become familiar with the concept of sintering a diamond layer onto a substrate to form a PDC insert. For more information concerning this process, the reader is directed to U.S. Pat. No. 3,745,623, issued to Wentorf Jr. et al., on Jul. 7, 1973.\nMany existing art PDCs exhibit durability problems in cutting through tough geologic formations, where the diamond working surface can experience high stress loads which may be transient in nature. Under such conditions, typical PDCs have a tendency to crack, spall, and break. Similarly, existing PDCs are relatively weak when placed under high loads from a variety of angles. These problems of existing PDCs are further exacerbated by the dynamic nature of both normal and torsional loading during the drilling process, during which the bit face moves into and out of contact with the uncut material forming the bottom of the well bore.\nFor optimal performance, the interface between the diamond layer and the tungsten carbide substrate must be capable of sustaining the high residual stresses that arise from the thermal expansion and bulk modulus mismatches between the two materials. These mismatches can create high residual stress at the interface as the materials are cooled from the high temperature and pressure process. Residual stress can be deleterious to the life of the PDC cutting elements, or inserts, during drilling operations, when high tensile stresses in the substrate or diamond layer may cause fracture, spalling, or complete delamination of the diamond layer from the substrate.\nDiamond is used as a drilling material primarily because of its extreme hardness and abrasion resistance. However, diamond also has a major drawback. Diamond, as a cutting material, has very poor toughness, that is, it is very brittle. Therefore, anything that contributes to further reducing the toughness of the diamond, substantially degrades its durability.\nA number of other approaches and applications of PDCs are well established in related art. The applicant includes the following references to related art patents for the reader's general familiarization with this technology.\nU.S. Pat. No. 4,109,737 describes a rotary drill bit for rock drilling comprising a plurality of cutting elements mounted by interference-fit in recesses in the crown of the drill bit.\nU.S. Pat. No. 4,604,106 reveals a composite polycrystalline diamond compact comprising at least one layer of diamond crystals and precemented carbide pieces which have been pressed under sufficient heat and pressure to create composite polycrystalline material wherein polycrystalline diamond and the precemented carbide pieces are interspersed in one another.\nU.S. Pat. No. 4,694,918 describes an insert that has a tungsten carbide body and at least two layers at the protruding drilling portion of the insert. The outermost layer contains polycrystalline diamond and the remaining layers adjacent to the polycrystalline diamond layer are transition layers containing a composite of diamond crystals and precemented tungsten carbide, the composite having a higher diamond crystal content adjacent to the polycrystalline diamond layer and a higher precemented tungsten carbide content adjacent to the tungsten carbide layer.\nU.S. Pat. No. 4,858,707 describes a diamond insert for a rotary drag bit consists of an insert stud body that forms a first base end and a second cutter end.\nU.S. Pat. No. 4,997,049 describes a tool insert having a cemented carbide substrate with a recess formed in one end of the substrate and having abrasive compacts located in the recesses and bonded to the substrate.\nU.S. Pat. No. 5,154,023 describes a process for polishing refractory materials, including natural and synthetic diamond, wherein the surfaces are successively softened to a predetermined depth by io implantation, followed by mechanical polishing.\nU.S. Pat. No. 5,154,245 relates to a rock bit insert of cemented carbide for percussive or rotary crushing rock drilling. The button insert is provided with one or more bodies of polycrystalline diamond in the surface produced by high pressure and high temperature in the diamond stable area. Each diamond body is completely surrounded by cemented carbide except the top surface.\nU.S. Pat. No. 5,217,081 relates to a rock bit insert of cemented carbide provided with one or more bodies or layers of diamond and/or cubic boron nitride produced at high pressure and high temperature in the diamond or cubic boron nitride stable area. The body of cemented carbide has a multi-structure containing eta-phase surrounded by a surface zone of cemented carbide free of eta-phase and having a low content of cobalt in the surface and a higher content of cobalt next to the eta-phase zone.\nU.S. Pat. No. 5,264,283 relates to buttons, inserts and bodies that comprise cemented carbide provided with bodies and/or layers of CVD- or PVD-fabricated diamond and then high pressure/high temperature treated in the diamond stable area.\nU.S. Pat. No. 5,304,342 describes a sintered product useful for abrasion- and impact-resistant tools and the like, comprising an iron-group metal binder and refractory metal carbide particles.\nU.S. Pat. No. 5,335,738 relates to a button of cemented carbide. The button is provided with a layer of diamond produced at high pressure and high temperature in the diamond stable area. The cemented carbide has a multi-phase structure having a core that contains eta-phase surrounded by a surface zone of cemented carbide free of eta-phase.\nU.S. Pat. No. 5,370,195 describes a drill bit having a means for connecting the bit to a drill string and a plurality of inserts at the other end for crushing the rock to be drilled, where the inserts have a cemented tungsten carbide body partially embedded in the drill bit and at least two layers at the protruding drilling portion of the insert. The outermost layer contains polycrystalline diamond and particles of carbide or carbonitride.\nU.S. Pat. No. 5,379,854 discloses a cutting element which has a metal carbide stud with a plurality of ridges formed in a reduced or full diameter hemispherical outer end portion of said metal carbide stud. The ridges extend outwardly beyond the outer end portion of the metal carbide stud. A layer of polycrystalline material, resistant to corrosive and abrasive materials, is disposed over the ridges and the outer end portion of the metal carbide stud to form a hemispherical cap.\nU.S. Pat. No. 5,447,208 describes a cutting element having a polished, low friction substantially planar cutting face with a surface finish roughness of 10 mu inch or less and preferably 0.5 mu inch or less.\nU.S. Pat. No. 5,544,713 discloses a cutting element with a metal carbide stud that has a conic tip formed with a reduced diameter hemispherical outer tip end portion of said metal carbide stud. A corrosive and abrasive resistant polycrystalline material layer is also disposed over the outer end portion of the metal carbide stud to form a cap, and an alternate conic form has a flat tip face. A chisel insert has a transecting edge and opposing flat faces, which chisel insert is also covered with a polycrystalline diamond compact layer.\nU.S. Pat. No. 5,624,068 describes buttons, inserts and bodies for rock drilling, rock cutting, metal cutting and wear part applications, where the buttons or inserts or bodies comprise cemented carbide provided with bodies and/or layers of CVD- or PVD-fabricated diamond and then HP/HT treated in a diamond stable area.\nU.S. Pat. No. 5,653,300 describes a superhard cutting element having a polished, low friction substantially planar cutting face with a surface finish roughness of 10 mu inch or less and preferably 0.5 mu inch or less. A chamfered cutting edge and side surface of the superhard material table of the same surface finish roughness are also disclosed.\nEach of the aforementioned patents and elements of related art is hereby incorporated by referenced in its entirety for the material disclosed therein."}
{"text": "Mirrors are typically used in either hand held or mounted form to assist in applying make-up, inspecting the face and the eyes, etc. It is especially desirable to use a mirror in close proximity and with a certain degree of magnification to assure proper execution of detailed procedures such as application of make-up, removing hair or treating blemishes.\nIt is known to provide a hand held or stand-mounted mirror which provides a normal sized mirror reflection on one side having a mirrored surface on the other side which provides a magnified reflection. However, the degree of magnification is limited by the physical construction of the mirror. As magnification increases in such a mirror, substantial distortion is encountered. Also, as one moves closer to the magnifying mirror surface, the degree of magnification decreases.\nIn U.S. Pat. No. 3,751,140, a combination of magnifying and non-magnifying mirrors are used for close-up viewing of the eyes and face. However, no means are provided for producing a higher degree of magnification for example, by the use of a lens. The invention concerns orienting the mirrors to provide images to the viewer from different angles of perspective.\nIn U.S. Pat. No. 2,817,998, a mirror camera has an objective formed by a concave mirror to which one or more correcting elements are added. The camera additionally incorporates means to prevent distortion of the mirror through impact or shock damage. The mirror camera utilizes a concave mirror with a meniscus shaped lens as a correcting element. There is no description of the utilization of a mirror with a magnifying lens for providing magnified viewing of a reflection for personal use.\nU.S. Pat. No. 2,584,829, discloses a hand mirror and supporting handle which incorporates a closure cap to protect the mirror faces."}
{"text": "In the last century reconstructive and cosmetic surgery has become a common practice. Specifically cosmetic breast surgery has been developed to allow reconstruction of a woman's breast that was affected by procedures such as mastectomy. Cosmetic breast surgery has also become available to amend the appearance of a woman's breast, for example by adding an implant to increase the size of the breast, to correct asymmetries, change shape and fix deformities.\nGenerally the implant is required to be able to provide a specific form and maintain the form for many years, preferably for the lifetime of the woman in which the implant is installed to prevent the need for additional invasive surgery. The implant is also required to have a specific feel preferably imitating the feel of a real breast. The implant also needs to be bio-durable such that it is not ruined by interaction with the human body and it needs to be bio-compatible so that the woman's health is not detrimentally affected by the implant even under extreme circumstances, for example the implant is required to be non toxic in case of leakage from the implant.\nThe standard implants used today comprise an outer shell typically formed from vulcanized silicone or polyurethane, and an inner content typically formed from a silicone gel or saline. The specific weight of the commonly used filling materials is generally between 0.95 to 1.15 (grams per centimeter cube). An average implant may weigh between 50 to 1000 grams, or even more. The weight of the implant is an addition, which is not negligible for a person.\nOver time breast implants are known to cause many problems, mostly related to the weight of the implant, for example: ptosis (i.e. sagging and deformity), breast tissue atrophy, prominence of the implant through breast tissue, back pain, and striae of the skin.\nTraditionally, the silicone gels used had silicone oils with short polymers in them that leached out through the shell over time. Current implants involve the use of a shell with barrier layers to achieve very low permeability, and using a cohesive gel as the filling material. The cohesiveness ensures that the filling material does not leak out into the body, even in case of rupture of the shell.\nAn additional characteristic to be considered in selection of the filling material is the resilience, elasticity and pliability of the implant, which provides it with a specific feeling when being sensed. Generally it is desirable to provide an implant which provides a specific shape and mimics the feel of real human tissue at the position of the implant. It is important that the implant maintain its form and feel for extended periods, to prevent the need for additional surgery.\nUS patent application publication no. 2004/0153151 to Gonzales dated Aug. 5, 2004 of which the disclosure is incorporated herein by reference describes a breast prosthesis from silicone that is formed as a trabecular body or micro-cell body in order to obtain a prosthesis of lower density.\nU.S. Pat. No. 4,380,569 to Shaw dated Apr. 19, 1983 of which the disclosure is incorporated herein by reference, describes a reduced weight breast prosthesis which is worn external to the human body or implanted into the human body. The breast prosthesis is comprised from a mixture of a silicone gel with glass micro-spheres.\nU.S. Pat. No. 5,902,335 to Snyder, Jr. dated May 11, 1999 of which the disclosure is incorporated herein by reference, describes a reduced weight breast prosthesis which is worn external to the human body, based on the U.S. Pat. No. 4,380,569 to Shaw. Snyder states that the use of glass micro-spheres as described by Shaw results in a stiff product that does not mimic the human breast as well as silicone gel alone. Snyder describes a breast prosthesis having two sections. A first outer section filled with silicone gel that mimics the human breast and a second inner section of reduced weight to reduce the weight of the prosthesis.\nU.S. Pat. No. 5,658,330,to Carlisle et al. dated Aug. 19, 1997 of which the disclosure is incorporated herein by reference, describes a molded silicone foam implant and method for making it."}
{"text": "The explosive growth of the Internet as a publication and interactive communication platform has created an electronic environment that is changing the way business is transacted. As the Internet becomes increasingly accessible around the world, online communications and business transactions increase exponentially.\nSeveral attempts have been made to facilitate online management of financial data using such network-based communications, namely to provide software packages residing on a computer and configured, for example, to acquire data from network-based financial transaction facilities, and to facilitate organization and management of such data over the network. However, such packages do not provide a satisfactory level of data access and/or detail. For example, many of the current software applications rely on preparation and display of transaction lists, such as credit card monthly statements and/or bank account statements, but do not provide a system to organize such transactions. Other software applications allow categorization of transactions, but require user input and data downloading prior to such user-defined operations. Thus, what is needed is a method and apparatus that facilitate real-time data retrieval, organization, and management over a network, such as the Internet."}
{"text": "The present invention relates to a head amplifier for use with a home VTR, more specifically, to a head amplifier suitable for integratin into an IC.\nA head amplifier is an amplifier for a home VTR which processes playback signals from a magnetic head. A head amplifier installed in an IC is disclosed, for example, in Japanese Laid-Open Publication No. 55-153108. The head amplifier of this publication has several disadvantages which will be detailed hereinbelow."}
{"text": "This invention relates to a filter with a valve-in-head construction. The filter is adapted for household and other uses for filtering fluid flowing through a pipeline. The above identified patents illustrate filters of this general type. Normally, such filters include an outer housing which is adapted to receive a replaceable filter cartridge. One end of the housing is closed, while the other end is adapted to receive a removable valve head which connects to the fluid line for selective fluid flow through the filter.\nVarious constructions of valve heads are shown in the above patents. Some of them, such as U.S. Pat. No. 3,777,889, include an annular \"stand pipe\" which telescopes into the core portion of the filter to locate it. Some also include inlets and outlets which are diametrically opposed, as in U.S. Pat. No. 4,271,020. Others, such as U.S. Pat. No. 3,853,761, utilize suitably actuated sliding spools which control the fluid flow within the valve. Some known devices also provide a bypass of portions of the fluid around the valve, even when it is in the \"off\" position.\nIt is an object of the present invention to provide a valve-in-head filter arrangement which is improved over the previously known structures.\nIn accordance with the various aspects of the invention, the device is for use with an outer housing adapted to receive a fluid filter cartridge therein. A valve head is provided which includes an annular stand pipe for locating an end of the filter. The stand pipe is cut away on its downstream side to reduce fluid back pressure. Furthermore, a passage is disposed between the head inlet and outlet, with the passage receiving a slideable plunger therein. An actuatable device connects through a rotary actuator to move the plunger. The actuator includes a stem mounted eccentrically to the handle axis, with the stem received in a recess in the plunger so that turning of the handle moves the plunger between \"on\" and \"off\" positions. The plunger includes a seal which continuously engages the passage wall in all plunger positions, and the construction provides full \"on\" and \"off\", with no fluid bypass in the \"off\" position."}
{"text": "The escalating demands for high density and performance associated with ultra large scale integration semiconductor devices require design features, such as gate lengths, below 100 nanometers (nm), high reliability and increased manufacturing throughput. The reduction of design features below 100 nm challenges the limitations of conventional methodology.\nFor example, when the gate length of conventional planar metal oxide semiconductor field effect transistors (MOSFETs) is scaled below 100 nm, problems associated with short channel effects, such as excessive leakage between the source and drain, become increasingly difficult to overcome. In addition, mobility degradation and a number of process issues also make it difficult to scale conventional MOSFETs to include increasingly smaller device features. New device structures are, therefore, being explored to improve FET performance and allow further device scaling.\nDouble-gate MOSFETs represent structures that have been considered as candidates for succeeding existing planar MOSFETs. In double-gate MOSFETs, two gates may be used to control short channel effects. A FinFET is a recent double-gate structure that exhibits good short channel behavior. A FinFET includes a channel formed in a vertical fin. The FinFET structure may be fabricated using layout and process techniques similar to those used for conventional planar MOSFETs."}
{"text": "The invention relates to a method of stripping metal sheets coated with molten material and in particular to stripping sheet or sheet material coated in a galvanising or vitreous-coating installation directly after the coated sheet has left a bath of molten coating material, both sides of the strip-like or web-like sheet being blasted with a thin curtain of gas. In addition the invention relates to apparatus for carrying out this method with two opposed gas nozzles which are arranged near and above the bath from which the sheet is drawn, and which form a thin curtain of gas for blasting both sides of the strip-shaped sheet or sheet-like material passing between them.\nIt is known, in the case of sheet and in particular strip-shaped sheet material which has been coated with molten coating material in a galvanising or vitreous-coating installation, to blast the surfaces of the sheet with a thin curtain of gas in order to strip off surplus coating material and prevent lumps or ribs of the coating material forming on the surfaces and thereby not only giving rise to unattractive surfaces, but also making the coating unnecessarily thick at different regions, representing an unnecessarily high consumption of coating material.\nIt is known to blast the surfaces of the coated sheet with air. However this leads to oxidation or other actions on the coating material, which again can result in flaws in the coated surfaces. Accordingly there has been a move towards using an inert gas such as nitrogen for blasting the surfaces of the coated sheet. It is true that this does satisfactorily eliminate the oxidation problem, but the gas consumption is relatively high, leading to undesirable increaases in cost."}
{"text": "Micro and nano structures including nanoparticle assembly with two and three dimensional periodicity can have potential applications in the areas of photonic crystals, chemical sensors, catalysts, and biotechnology. Patterned surfaces can be used as hard templates to assist the self assembly of not only relatively simple clusters but also complex and unique crystallization structures. Soft polymer templates have been used for directed self assembly of particle arrays on flat substrates. Binary colloidal crystals have been fabricated using two different sizes of colloidal particles. Further, micro and nano particles have been used as templates for the preparation of porous metallic nanostructures and monodisperse colloidal crystals. Even though nanochannel structures for nanofluidic applications have been fabricated using thermal oxidation or nanoimprint, there is a need for a simple and inexpensive approach for the fabrication of enclosed channels formed of nanoparticles.\nThus, there is need to solve these and other problems of the prior art and provide a simple method for the fabrication of nanochannel structures."}
{"text": "There is a known HMD that is worn around a user's head and can present the user with an image in a virtual space with the aid of a display or any other device disposed in front of the user. In particular, the HMD disclosed in Patent Document 1 can display a 360-degree panoramic image in a three-dimensional virtual space. Such an HMD typically includes a variety of sensors (an acceleration sensor and an angular velocity sensor, for example) and measures data on the attitude of an HMD main body. In particular, the direction of the sight line of the eyes looking at the panoramic image can be changed in accordance with information on the angle of rotation of the head. That is, when the user who wears the HMD rotates his/her own head, the direction of the sight line of the eyes looking at the 360-degree panoramic image is changed accordingly, whereby the user's immersive sensation in the video world is enhanced for improvement in entertainment quality (see paragraphs [0004], [0005], [0006], and [Abstract] of Patent Document 1).\nPatent Document 1 further discloses that the user's gesture in the form of motion of the head is related in advance to operation (display screen switching, for example) on the screen of the HMD, and when the user's gesture is identified on the basis of the inclination of the head or the acceleration thereof, screen operation related to the gesture is performed. Inconvenience in controller operation resulting from the fact that the user who wears the HMD is unable to see his/her own hands and therearound is therefore solved (see paragraphs [0008] to [0009] in Patent Document 1)."}
{"text": "DE 10 2012 010 757 A1 discloses an illuminating device for a vehicle. The illuminating device includes an eye-tracking system for detecting an eye position and its viewing direction. The illuminating device is divided into several predefined switch-on ranges, and light is emitted only into that predefined switch-on range which corresponds to the user's field of vision. No light is emitted as soon as the user's field of vision does not correspond to the respective predefined switch-on range.\nThe disadvantage, however, consists in that drivers of a motor vehicle during darkness frequently orientate themselves using the lighting of motor vehicle instruments. If the instruments are always completely switched-off as soon as the driver ceases to look in their direction, the driver loses his orientation more easily, which may make operation of the motor vehicle more difficult.\nThis leads to the requirement to further develop a method to address this disadvantage in such a way that it becomes easier for the driver to orientate himself in a dark motor vehicle without being irritated by the lighting of at least one motor vehicle instrument."}
{"text": "(i) Field of the Invention\nThis invention relates to a waterproof seal ring to be fitted between the contacting surfaces of a master cylinder body and a power booster housing.\n(ii) Description of the Prior Art\nConventional waterproof seal rings of this character have a recess formed at the lower end of the outer periphery, through which air is allowed to flow into and out of the master cylinder as its piston reciprocates. When rainwater or the like gains entrance between the contacting surfaces of the master cylinder body and the power boosting housing, it will flow down around the outer periphery of the seal ring. The intruding liquid can then be entrained by air that is drawn by suction into the master cylinder, eventually to corrode its inner wall surface. This is particularly true with the arrangements in which the master cylinders are installed with their axes increasingly raised frontwards, for example, to meet the recent design requirements of the braking systems for the front-engine front-wheel-drive automobiles."}
{"text": "This invention relates to a pneumatic tire having an asymmetric tread pattern. More particularly, it relates to a pneumatic tire capable of improving driving performance on a wet road surface and reducing noise during driving while keeping excellent driving performance on a dry road surface.\nA pneumatic tire having an asymmetric tread pattern wherein ribs continuing in a tire circumferential direction are disposed at a shoulder end portion on a tread surface on the outer side of a vehicle when the tire is fitted to the vehicle, and block lines each comprising a plurality of blocks are disposed at a shoulder end portion on the inner side of the vehicle exhibits excellent driving performance on a dry road surface because its outer shoulder end portion is constituted by ribs having high rigidity. It is known generally that such a pneumatic tire exhibits excellent wear resistance even under a critical state such as in circuit driving because the shoulder end portion on the outer side of the vehicle comprises the ribs having high rigidity and the tire does not have relatively sharp block edges.\nHowever, though this pneumatic tire has no problem of draining property on a wet road surface having a relatively small depth of water, it involves inferior drainability when the depth of water is great because the shoulder portion on the outer side of the tire with respect to the vehicle consists of the ribs. Various attempts have been made so far to solve these problems by, for example, increasing the void volume of the tread, or by employing an asymmetric profile as disclosed in Japanese patent application Kokai publications No. 57-147901 and No. 61-98601. However, these methods have not been entirely satisfactory because they invite the drop of other items of tire performance.\nAlso, lately it has been increasingly strongly demanded from the viewpoint of the environmental integrity that noise generation occurring when vehicles are run should be suppressed, and the above described tires having remarkable dry performance are not excluded in this respect. Tires the tread pattern of which is asymetrical and comprises a rib at the shoulder end on the outer side in the condition in which they are mounted to wheels of a vehicle have a feature such that they do not have a continuous lug groove in the outer rib in contrast to such tires of which the tread pattern comprises a block-based pattern over a whole tread area, and basically such feature ought to be advantageously influential upon supression of the noise generation. However, a pair of tires, one for a left and the other for a right wheels of a vehicle, are produced by a same molding die, so that in mounting such tires to the vehicle, either of the two tires has to be reversed in the mounting-position relative to the other so that the outer rib of each tire is on the outer side of the vehicle. If the tire mounting is made in that way, then a basic direction of inclination of sub-grooves extending in the widthwise direction of the tires falls to be identical in connection with the left tire and the right tire on the vehicle, whereby a difference in the block rigidity is produced between the left tire and the right tire, possibly resulting in a large difference in the noise generation between the left side and the right side of the vehicle."}
{"text": "1. Field of the Invention\nThe present invention relates to microfluidic devices and methods.\n2. Description of Related Art\nTraditional methods for crystal growth and crystallization are highly labor intensive and require significant quantities of material to evaluate and optimize crystal growth conditions. Examples of these methods include the free interface diffusion method (Salemme, F. R. (1972) Arch. Biochem. Biophys. 151:533-539), vapor diffusion in the hanging or sitting drop method (McPherson, A. (1982) Preparation and Analysis of Protein Crystals, John Wiley and Son, New York, pp 82-127), and liquid dialysis (Bailey, K. (1940) Nature 145:934-935).\nPresently, the hanging drop method is the most commonly used method for growing macromolecular crystals from solution, especially for protein crystals. Generally, a droplet containing a protein solution is spotted on a cover slip and suspended in a sealed chamber that contains a reservoir with a higher concentration of precipitating agent. Over time, the solution in the droplet equilibrates with the reservoir by diffusing water vapor from the droplet, thereby slowly increasing the concentration of the protein and precipitating agent within the droplet, which in turn results in precipitation or crystallization of the protein.\nThe process of growing crystals with high diffraction quality is time-consuming and involves trial-and-error experiment on multiple solution variables such as pH, temperature, ionic strength, and specific concentrations of salts, organic additives, and detergents. In addition, the amount of highly purified protein is usually limited, multi-dimensional trials on these solution conditions are unrealistic, labor-intensive and costly.\nA few automated crystallization systems have been developed based on the hanging drop methods, for example Cox, M. J. and Weber, P. C. (1987) J. Appl. Cryst. 20:366; and Ward, K. B. et al. (1988) J. Crystal Growth 90:325-339. More recently, systems for crystallizing proteins in submicroliter drop volumes have been described including those described in PCT Publication Nos. WO00/078445 and WO00/060345.\nExisting crystallization, such as hanging drop, sitting drop, dialysis and other vapor diffusion methods have the limitation that the material for analysis and the crystallization medium are exposed to the environment for some time. As the volumes of materials decrease, the ratio of surface area to volume ratio varies as the inverse of the radius of the drop. This causes smaller volumes to be more susceptible to evaporation during the initial creation of the correct mixture and during the initial period after the volume has been set up. Typical hanging drop plates can have air volumes of 1.5 milliliters compared to a sample drop size of 3-10 microliters. Moreover, typical methods expose the sample drop to the environment for a duration of seconds to minutes. Small variability in the rate that samples are made can cause significant variations in the production of crystals. Small variations external environment also can cause significant variations in the production of crystals even if the rate that the samples are made is unchanged. Prior methods fail to reduce the problems of convection currents under 1 g such as those described in U.S. Pat. No. 4,886,646, without the large expenditure of resources or in methods that complicate crystal analysis."}
{"text": "Tape drives that operate bi-directionally are well-known. The tape media is generally enclosed into a single reel cartridge and the tape is transported around the tape path onto a tape reel thereby placing the tape media in contact with the transducer. The transducer in present day tape drives comprise a separate read and write element that cover each track of the data on the tape media. It is known, for instance, that a one-half inch wide magnetic tape can include 18 tracks. A higher number of tracks are contemplated and it is proposed that a one-half inch wide magnetic tape includes 36 tracks. To be able to read 36 tracks from the tape, an interleaved data transducing head is proposed.\nAn interleaved read/write head is disclosed in U.S. Pat. No. 4,685,005 to Fields, and assigned to the assignee of the present invention. In the magnetic head disclosed in that patent, the read and write gaps of each module of the magnetic head are alternately spaced across the width of the tape, such that the write gaps of one module are aligned with the read gaps of the other module. When one module is selected for writing, as a function of the direction of the tape movement, the other module is selected to read-after-write check the data written by the write element of the first module. One module writes odd track data during one direction of the tape movement and reads even track data during the opposite direction of the tape movement. The problem arises in that both read and write elements are located within one module with the gaps of each element aligned along the same line.\nThe read element of the modern day tape head is a magneto-resistive transducer and is formed of thin film layers deposited through a standard thin film deposition procedure. The write transducer, however, has its pole pieces formed from two blocks of magnetic ferrite magnetically connected together at the back gap and a transducing gap at the front face. A single block of magnetic ferrite operates as a closure structure after all of the elements of the read transducer and the write conductors are deposited onto the substrate magnetic ferrite block. Having different layers of material deposited at different locations into the side-by-side elements creates a leveling problem when the magnetic ferrite closure block is to be placed over the plurality of read and write elements. More layers form a greater thickness in the read element and therefore the write element generally lacks a supporting structure that can cause the magnetic ferrite closure piece to bend under stress when affixed to the completed transducers. In any event, the write head gap must not be narrower than the read head gap to properly read the written bits of information onto the magnetic media. Also, with the closure piece bent under stress at the time of manufacture, the stress may relieve in time, resulting in an unreliable gap.\nIt is, therefore, an object of the present invention to provide an enhanced magnetic recording drive for multi-track operation.\nYet another object of the present invention is to provide a magnetic drive that uses an enhanced interleaved head to read and write magnetic transitions forming data onto magnetic recording media.\nAnother object of the present invention, therefore, is to provide an enhanced transducer assembly and an enhanced method for making the transducer assembly."}
{"text": "1. Field of the Invention\nThis invention relates a parts welding device, wherein a part, such as a bolt, is held by a spring chuck mean and the part is welded to the other part, such as a steel plate.\n2. Description of Related Art\nAs a method of welding a bolt where a flange is integrally formed on the shaft to a steel plate, stud welding is known. In this method, a collet-type chuck with a spring structure is arranged on the head of a feed pipe for the bolt and the shaft of the bolt moved to the chuck is held by the spring structure, wherein the flange is exposed on the head side. The flange is brought closer to the steel plate, and welding current is made to flow between the flange and the steel plate to generate arc, which fuses the flange and steel plate together by metal fusion. Then, pressure welding follows, completing the welding process.\nThe above method is to hold the bolt after stopping the approaching bolt by the collet-type chuck. This method, however, raises the problem that the holding position of the bolt is not determined in a constant manner. In such a situation, the gap between the flange and steel plate becomes not constant, causing the unevenness of welding quality. Further, for automating such a welding system, feeding and holding the bolt in a secure manner upon welding is imperative, but the above method cannot provide an operation enabling the feeding and holding to be made in a satisfactory manner. Besides, in the case of a nut with a flange integrally formed on the cylindrical part of the nut, on the internal surface of which a female screw is formed, a similar problem occurs when the flange is welded to the other part.\nThis invention is a parts welding device provided to solve the problem mentioned above. In the invented system, a part feed apparatus and a welding apparatus are mounted on a board, a supporting member is fitted on a supporting shaft fixed on the board in a rotary manner, a welding unit is fitted into the supporting member, and a receiving hole for a part fed by the part feed apparatus is formed on the head of the supporting shaft, wherein the head tip of the welding unit is arranged near the opening of the receiving hole.\nThe part fed by the part feed apparatus is inserted into the receiving hole and held there, and the part comes in contact with the other part, then both parts are welded together by the welding unit. The supporting shaft has the opening as the receiving hole for the part and the supporting member is fitted on the supporting shaft in a rotary manner, that is, the supporting member equipped with the welding unit is fitted on the supporting shaft in a rotary manner. Therefore, the operations, such as holding the part in the receiving hole and fitting or rotation of the welding unit, are arranged in a compact manner while the supporting shaft is set as the central member for the operations, providing a welding system having a constitutional advantage. In addition, when the part is fed to the receiving hole, the part is inserted in the hole from the outside, so that it can be held with the supporting shaft in a smooth manner.\nIt is applicable that the above supporting member is rotated by a rotation mean so that the head tip is moved to a prescribed position. As the supporting member rotates, the welding unit fitted into the supporting member rotates together, allowing the head tip to move freely to a required welding spot. Therefore, when a circular flange is welded to a steel plate, local welding spots can be arranged apart at 120 or 180 degree without fail.\nWhile the above supporting shaft is arranged in the direction almost perpendicular to the board and the above part feed apparatus inserts the part retained on the feed rod of the feed unit into the above receiving hole from the head side of the above supporting shaft, a retaining mean for the part may be arranged in the above receiving hole. As described above, the part feeding behavior is to inserting the part into the receiving hole form the head side of the supporting shaft, wherein part holding can be made by the most simple method of inserting the part from the outside. As a result, the dislocation of the holding position of the part never occurs, making the relative position between the part and the other part constant, thus a good welding quality can be obtained. Further, arranging the part holding mean in the receiving hole keeps the relative position between the receiving hole and the part constant, maintaining a correct contact state upon making the part come in contact with the other part, thus works effectively to improve welding quality.\nWhile the feed rod of the above part feed apparatus is made to move back and forth in an inclined state against the above supporting shaft, a moving-back-and-forth mean may be arranged for moving back and forth the whole body of the above feed rod in the same direction of the axis of the part retained on the feed rod. The part retained on the head of the inclined feed rod stops in the position where the axes of the part and receiving hole are aligned. Then, the whole body of the feed rod moved in the axial direction of the receiving hole and the part, allowing the shaft of the part to be inserted smoothly into the receiving hole. As described here, by arranging a proper position relation between the feed rod and supporting shaft and moving the whole body of the feed rod, the above smooth operation can be achieved, so that a high credibility as a welding system can be maintained. In order to secure the above coaxial state between the shaft and receiving hole, in this system, the axis of the supporting shaft, the axis of the shaft held by a supporting pipe, and the stroke direction of the piston rod of an air cylinder are all set in parallel. Thus, when the whole body of the feed rod is moved, the shaft is inserted smoothly in the receiving hole.\nThe above board can be fitted on the head of a robot apparatus. By operating the robot apparatus, the board is moved freely to any directions and the direction of the board upon its stopping can be selected freely, the parts welding can be made without any hamper even if the other part has a complicated form or the welding point is in a complicated area. In addition to such a free moves of the board, the particular arrangement made for the supporting shaft and feed rod gives more advantages to this welding system, further enhancing the utility of the system.\nIn another application, the above board is fixed to a stationary member and the other part to which the part is welded is held with the robot apparatus. Such a constitution is reverse to the one mentioned above, but secures the same advantages obtained by the one mentioned above. In the reversed constitution, the other part is held with the robot to move freely, while the part feed apparatus and the welding apparatus are kept stationary."}
{"text": "1. Field of the Invention\nThe present invention relates to the use of acoustic waves for charging batteries in a battery charger or in a portable telephone.\n2. Description of the Related Art\nHome audio systems use audio speakers to convert audio signal into acoustic waves that can be heard by a user. The audio signals are analog electronic signals that represent the time-variation of the desired acoustic waves. They generally represent music or voice signals with frequency components typically in the range of 50 Hz-20 kHz. Typically, the audio signals are received by an amplifier from a radio receiver, compact-disk player, cassette player, or other audio source, amplified to appropriate voltage and current levels, and provided to one or more audio speakers to generate the corresponding acoustic waves.\nThus, audio speakers are energy transducers that convert electrical energy into acoustic energy. The transduction from acoustic energy to electrical energy is typically performed in the speakers by providing the audio signal to a solenoid to generate a time-varying magnetic field. The time-varying magnetic field then displaces a magnetic element coupled to a speaker diaphragm, causing the speaker diaphragm to vibrate with a displacement indicated by the audio signal. With the audio speaker placed in a listening area, the vibrating speaker diaphragm generates the desired acoustic waves in the air of the listening area.\nRechargeable batteries and battery chargers for rechargeable batteries are well known in the art. The rechargeable batteries are based on a variety of electrochemical reactions, and are generally recharged by applying an appropriate potential across the rechargeable battery to run a current through the rechargeable battery. Some battery chargers monitor the energy stored in a rechargeable battery and adjust the applied potential to optimize the charging process. A main feature of battery chargers is their use of some electrical power source, such as home 120 V AC electrical power source, to drive a current through a rechargeable battery.\nIt would be useful to have a battery charger that does not draw current from other batteries or from a home 120 V AC electrical power source, but which instead uses an otherwise unused source of power, such as the acoustic energy generated in a room by people in the room performing typical activities such as talking, walking, and moving."}
{"text": "Generally, for health maintenance (improvement of lifestyle-related diseases) or enhancement of athletic abilities, it is valid to decrease body fats and promote carbohydrate metabolism by aerobic exercises. That is to say, if we incorporate glucoses which are energy sources of muscles into the muscles and burn the glucoses, we consume surplus glucoses, as a result of improvement of hyperglycemia and hyperinsulinemia, incorporating the glucoses and burning the glucoses contribute the improvement of lifestyle-related diseases (diabetes, obesity and hyperlipemia and so on).\nThe incorporating the surplus glucoses into the muscles is caused by muscle contractions, if we increase incorporated quantity into the muscles, it becomes possible to promote the carbohydrate metabolism.\nIn order to do the carbohydrate metabolism effectively by the muscle contractions, it is desirable to cause the muscle contractions of large volume muscles, we think that it is valid to extend and bend a hip joint and contract a big muscle group from a trunk to femoral regions.\nThere are running or walking et cetera as exercises contracting the big muscle group, and in the capacity of a device to imitate those exercises, treadmill is generally known.\nThe treadmill is an exercise assisting device to be capable of adapting rotational velocity and inclination angle of a moving belt which is a walking surface or running surface. And, as an aim of a training of a whole body endurance, running exercise or walking exercise on the treadmill are often done.\nOn the other hand, in the capacity of an exercise assisting device which is different from the treadmill and which intend a simulation of running or walking, it is known that an axle moving type bicycle ergometer which is disclosed in Japanese Patent Application Laid Open No. 2002-78817.\nThe ergometer comprises a pair of rotational axes which are given each rotational force by a pair of right and left pedals and a pair of right and left arms, a pair of tables which supports each axis of the pair of rotational axes in a horizontal direction, a pair of guide parts which supports each table of the pair of tables capably of reciprocating in a direction orthogonal to the pair of rotational axes, a pair of moving mechanisms which move the pair of tables on the pair of guide parts based on rotational angle of the rotational axis and a braking mechanism which do a breaking of each rotational axis of the pair of rotational axes (refer to Japanese Patent Application Laid Open No. 2002-78817)."}
{"text": "As for the engine torque control of an internal combustion engine controller for automobiles, so-called “torque base (torque demand) type engine control” in which target engine torque is computed from the angle of an accelerator and the speed of an engine and throttle control, fuel control and ignition control are carried out to achieve the target engine torque and the target air-fuel ratio has recently been implemented.\nThe torque base (torque demand) type engine control of the internal combustion engine has advantages that a torque difference at the time of switching between uniform charge combustion and stratified charge combustion in a stratified charge lean combustion system can be reduced and that traction control and engine torque demanded from an external device such as an auto-cruise or AT can be processed smoothly by adding an interface for external demand torque to a logic for computing the above target engine torque.\nThe above torque base type engine control has an advantage that torque control can be carried out while the target air-fuel ratio is maintained as torque control is basically carried out by controlling the quantity of air sucked by an electrically controlled throttle. However, it has a problem that its response for achieving desired torque is low due to a phenomenon that the supply of intake air into cylinders is delayed. To cope with this problem, when high-speed response is desired for traction control or VDC (vehicle dynamics control), for example, other torque control means is used in combination to improve the torque response. As one example of this, there is known technology making use of a fuel cut or ignition retard when torque is reduced (deceleration).\nAnother technology concerning the torque response is disclosed by JP-A H11-72033, for example. In this technology, when high-speed torque response is desired in a stratified charge lean combustion system, the torque response is improved by correcting the ignition time at the time of uniform charge combustion and the air-fuel ratio at the time of stratified charge combustion. In this technology, torque correction is carried out by correcting the ignition time as the purification efficiency of a three-way catalyst is reduced to deteriorate exhaust gas when the air-fuel ratio is corrected at the time of uniform charge combustion, and torque correction is carried out by correcting the air-fuel ratio in the case of stratified charge combustion as the variable range of ignition time is small.\nMeanwhile, as for the recent exhaust gas purification control of an automobile internal combustion engine, there is generally known technology for improving the exhaust purification ratio with a three-way catalyst by carrying out air-fuel ratio feed-back control, using a detection signal from an O2 sensor installed in the exhaust pipe, so that the air-fuel ratio becomes a value close to the theoretical air-fuel ratio.\nHowever, the three-way catalyst has an O2 storage effect (effect of storing oxygen in a catalyst) and the function of reacting with an exhaust component in the catalyst so that the stored O2 cancels a shift from the theoretical air-fuel ratio in exhaust from the engine. Therefore, when air-fuel ratio feed-back control is carried out by using information only from the O2 sensor without considering the exhaust purification function of the stored O2, the correction quantity of fuel becomes inappropriate and over-correction occurs, thereby deteriorating the exhaust gas. To cope with this problem, for example, technology disclosed by JP-A H2-230935 is to prevent the deterioration of the purification of exhaust gas by adjusting the amount of air-fuel ratio feed-back control based on an estimated value obtained by O2 storage quantity estimation means for estimating the storage quantity of O2 in a three-way catalyst provided in an internal combustion engine controller.\nFor the torque base type control of the conventional internal combustion engine controller for automobiles, technology for improving the torque response by using other high-speed torque control means is used to compensate for a delay in the supply of intake air when high-speed torque response is desired. In a uniform charge stoichiometric combustion system which is operated at a value close to the theoretical air-fuel ratio, there has not been proposed appropriate torque assist means when high-speed torque increase is demanded during non-idling.\nOne of the reasons for this is that torque increase is impossible at the time of non-idling by further advancing the ignition time as a value (MBT) at which the ignition time is advanced to enable the generation of the maximum torque is generally set as a standard ignition time even when torque assist is tried by changing the ignition time.\nAnother reason is that when the air-fuel ratio is simply made rich, the purification of exhaust gas may be deteriorated (specifically, increases in the quantities of CO or HC) as the air-fuel ratio must be maintained at a value close to the theoretical air-fuel ratio due to the exhaust gas purification properties of the three-way catalyst used in a uniform charge stoichiometric combustion system though it is commonly known that the torque can be increased by making the air-fuel ratio rich (for example, power air-fuel ratio of about 12) as for the technology of assisting torque by the air-fuel ratio control of an internal combustion engine.\nIt is an object of the present invention which has been made in view of the above problems to provide an internal combustion engine controller for automobiles which can realize torque increasing performance and exhaust gas purification performance in a well-balanced manner even when a high-speed torque increase is demanded in uniform charge stoichiometric combustion which is operated at a value close to a theoretical air-fuel ratio in the torque base type engine control of an internal combustion engine."}
{"text": "Metallurgical plants are plants for processing metal ore, wherein the central element of such a metallurgical plant is a blast furnace. These metallurgical plants have been known for a long time. A blast furnace is fed with raw materials which comprise metal ore, additives and heating material. Usually coal or coke is used as a heating material, wherein coal and coke produce heat by burning in the presence of air on the one hand and wherein coal and coke also function as reduction agent for the metal ore, as the metal ore is basically comprised of metal oxide. When reducing metal ore in a blast furnace, various gases are produced, which collectively are known as a furnace gas or flue gas. Said furnace gas usually contains a substantial amount of carbon dioxide (CO2). Carbon dioxide is a greenhouse gas, and during recent years more and more effort has been made to avoid or convert greenhouse gases, as these greenhouse gases are regarded as detrimental for the climate.\nIn the field of metal production, it is a general aim to use as few raw materials and heating materials as possible, as these materials are expensive and it is expensive to transport these materials. Much effort has been made to reduce the amount of coke/coal used in production. One approach was blowing coal dust into the blast furnace, and another approach was producing carbon monoxide as a reduction gas, either in the blast furnace itself or in a separate gasification reactor outside the blast furnace. From EP 09318401 A1 it is known to blow a portion of the carbon required for reducing the metal ore into the blast furnace in form of a substitute reduction material. In this sense, e.g. natural gas, heavy oil, fine coal and similar material having a high carbon content may be used as a substitute reduction material. These materials may be directly blown into the blast furnace shaft or may be gasified outside of the blast furnace shaft in a separate gasification reactor so as to form a reduction gas. Subsequently, such a reaction gas may be directed into the blast furnace shaft. The method known from EP 09318401 A1 may provide a possibility to reduce the consumed amount of coal or coke and may also provide a possibility to use materials, which are difficult to process as a substitute reduction material, however, the problem of high CO2 production in metal production has not been solved.\nThe prior art discloses methods wherein furnace gas or a particular component thereof is directed out of the blast furnace shaft and, after being processed in a CO2 converter, is re-directed into the blast furnace shaft. EP 2 543 743 A1 discloses a method wherein furnace gas is directed out of the blast furnace shaft and is directed to a separation device in which CO and CO2 are separated. Only the separated CO2 is subjected to reforming in a CO2 converter. Reforming produces mainly CO and H2O, wherein H2O is separated and CO is directed into the blast furnace shaft. WO 2011/087036 A1 also discloses a method wherein first furnace gas is directed to a separation device in which CO and CO2 are separated. In a CO2 converter, the CO2 is converted into O2 and CO. The CO from said conversion and the previously separated CO are jointly directed into the blast furnace shaft. U.S. Pat. No. 3,909,446 A discloses a method wherein furnace gas from a blast furnace shaft is mixed with coke oven gas in a CO2 converter. Thus, a gas mixture comprising CO and H2 is produced which is ro-directed into the blast furnace shaft. WO 2010/049536 A1 describes a similar method wherein also carbon containing particles are re-directed into the blast furnace shaft. U.S. Pat. No. 2,598,735 A discloses a method wherein furnace gas from a blast furnace shaft is mixed with carbon/coal and oxygen in a gas generator. A portion of the carbon is burnt in presence of the oxygen, and another portion of the carbon reduces the CO2 from the furnace gas and the CO2 from the burnt carbon to CO. Said CO is re-directed into the blast furnace shaft as a reducing agent. None of these documents discloses a method wherein further processing of a portion of the converted CO is carried out."}
{"text": "This invention relates to water purification systems and in particular to a control for a system of the type that purifies water by converting it into steam and condensing the steam back into water.\nWater distillation systems are well-known in the art and include a boiling tank having an inlet for admitting untreated feed water to the tank and heating means for heating the untreated water to produce steam. As steam is generated by the boiling tank, the minerals and other impurities in the residue water increase in concentration. Periodically, it is necessary to remove the residue by draining. In a water distillation unit designed for domestic home use, the drain commonly discharges the residue water into the house drain. With the advent of polyvinylchloride plastic pipe for plumbing systems, it is necessary to reduce the temperature of the residue water before discharging it to the drain. Accordingly, one object of the present invention is to provide an efficient and inexpensive means for reducing the temperature of the residue water before it is discharged to drain.\nEven if the residue is periodically discharged to drain, mineral scale will be deposited on inner surfaces of the tank and outer surfaces of the heater. Accordingly, it is necessary to periodically shut down the unit and descale the interior of the boiling tank by filling the tank with water having descaling chemicals added thereto and, after a soaking period, draining the cleaning liquid and rinsing the tank. It is an object of the present invention to provide a control having a cleaning mode that minimizes the complexity of the procedure performed by the user and that adds very little to the cost of the control.\nThe primary purpose of a water distillation unit is to produce distilled water and accumulate this water in a holding tank until needed for use. Consumers expect that such water will be of very high quality. If the water in the holding tank is contaminated by the condition of the holding tank then water quality may be diminished. Accordingly, it is known to disable the condensing portion of the water distillation apparatus in order to provide steam from the boiling tank directly to the holding tank in order to destroy microorganisms that may be established in the tank. The difficulty with known units having such a steam cleaning cycle is that, if the user forgets to reactivate the condensing means, then the steam will eventually condense within the holding tank providing a buildup of very hot water in the holding tank. If the holding tank outlet pump is operated by the user after inadvertently neglecting to discontinue the steam cleaning mode, it is possible that the hot water may cause damage to the holding tank outlet pump. Accordingly, it is an object of the present invention to provide an inexpensive control for protecting a water distillation unit against the failure of the user to discontinue the steam cleaning mode.\nAnother difficulty with water distillation units is that the addition of relatively cold untreated feed water to the boiling tank tends to rapidly lower the temperature of the interior of the tank. This lowering of temperature creates a vacuum within the boiling tank which tends to draw water already condensed in the condenser and steam previously produced back into the boiling tank where it will condense back into water. Therefore, the addition of feed water creates inefficiencies in the operation of the unit which results in increased energy consumption per unit of water produced. Accordingly, it is an object of the present invention to provide a control for a water distillation unit which carefully controls the frequency and amount of feed water admission to the boiling tank in order to reduce the detrimental affects described above."}
{"text": "1. Field of the Invention\nThe present invention relates to a circuit forming method and a circuit connection structure in an electrical connection box for use in a vehicle, for example.\n2. Description of the Related Art\nIn an electrical connection box of the sort mentioned above, an electronic circuit board 104 provided with electronic parts 103 between a lower case portion 101 and an upper case portion 102 constituting the electrical connection box, and an insulating board 106 provided with bus-bars 105 are generally housed in a laminated form as shown in FIG. 5.\nThe electrical connection of the electronic parts 103 to the bus-bars 105 is accomplished by fitting a connector housing 107 retaining female terminals 108 to an edge portion of an electronic circuit board 104 and stacking the electronic circuit board 104 in this state on top of the insulating board 106 thereby to fit bus-bar tabs (male terminals) 105a with the bus-bars 105 in an upright condition into the respective female terminals of the connector housing 107. An insulating spacer 109 is provided between the electronic circuit board 104 and the insulating board 106 to avoid unnecessary electrical connection.\nHowever, there have been the following problems concerning the connection structure of the bus-bar in the conventional electrical connection box.\n(1) When a wiring pattern (not shown) is formed for a number of electronic parts provided on the electronic circuit board 104, it is needed to form the wiring pattern up to the edge portion where the female terminals 108 are provided. Moreover, as a space for use in fitting up the female terminals 108 becomes necessary and this results in a large-sized electronic circuit board 104, thus increasing the cost.\n(2) The edge portion of the electronic circuit board 104 requires the female terminals 108 and the connector housing 107 so as to electrically connect the electronic parts 103 to the respective bus-bars 105, so that the number of parts tends to increase.\n(3) Moreover, because the bus-bar tab 105a connected to the female terminal 108 is formed by processing the edge portion of the bus-bar 105, any intermediate part of the bus-bar 105 cannot be connected to the female terminal 108 and when one bus-bar is connected to two or more connecting parts, for example, it is needed to use the same number of bus-bars as the number of connecting parts. Consequently, the utilization efficiency of the space for used in arranging the bus-bars has been poor.\nAn object of the present invention is to solve the foregoing conventional technical problems and to provide a circuit forming method and a circuit connection structure that are capable of making smaller an electronic circuit board as well as reducing the number of parts.\nThere is provided a circuit forming method in an electrical connection box according to the invention with an electronic circuit board and an insulating board being stacked up and housed therein, comprising the steps of providing not only a through-hole in the electronic circuit board but also lead wires or electric wires of electronic parts in a manner straddling the through-hole, providing the bus-bars on the insulating board so that the upright pressure-contact knife-edges of the respective bus-bars may be arranged in a position corresponding to the through-hole of the insulating board, and stacking up the electronic circuit board and the insulating board by moving both the boards close to each other in such a condition that the pressure-contact knife-edges are passing through the through-hole, and pressure-welding the lead wires or electric wires to the respective pressure-contact knife-edges.\nThere is provided a circuit connection structure comprising: an electronic circuit board having a through-hole, the electronic circuit board being provided with lead wires or electric wires of electronic parts in a manner straddling the through-hole, and an insulating board provided with bus-bars having upright pressure-contact knife-edges electrically connected to the lead wires or electric wires, wherein the electronic circuit board and the insulating board are stacked up by relatively moving both the boards close to each other in such a condition that the pressure-contact knife-edges are passing through the through-hole; and the lead wires or electric wires are pressure-welded to the respective pressure-contact knife-edges.\nAs the lead wires are pressure-welded to the respective pressure-contact knife-edges passed through the through-hole provided in the electronic circuit board in this circuit forming method and the circuit connection structure, the space for use in fitting the female terminals to the edge portion of the board in the prior art can be dispensed with, which results in making it possible to reduce the size of the electronic circuit board. Moreover, the pressure-contact knife-edges are directly pressure-welded to the respective lead wires, so that female terminals and the connector housing necessitated in the prior art can be omitted, which results in reducing the number of parts.\nIn the circuit connection structure according to the invention, the pressure-contact knife-edges are formed separately from the respective bus-bars, and the pressure-contact knife-edges may be fitted to the respective bus-bars in such a condition that the former has been electrically connected to the latter.\nWith this arrangement, as any number of pressure-contact knife-edges can freely be provided to the end portion or mid position of the bus-bar, the utilization efficiency of the space for use in disposing the bus-bars becomes improvable thereby."}
{"text": "Liquid phase separations (eg. liquid chromatography and electrophoretic separations) have long been used as investigative tools by scientists and researchers seeking to identify the structure of molecules, particularly peptides (as used herein the term “peptides” refers to polymers having more than one amino acid, and includes, without limitation, dipeptides, tripeptides, oligopeptides, and polypeptides. The term “protein” refers to molecules containing one or more polypeptide chains).\nProteomics involves the broad and systematic analysis of proteins, which includes their identification, quantification, and ultimately the attribution of one or more biological functions. Proteomic analyses are challenging due to the high complexity and dynamic range of protein abundances. The industrialisation of biology requires that the systematic analysis of expressed proteins be conducted in a high-throughput manner and with high sensitivity, further increasing the challenge. Recent technological advances in instrumentation, bio-informatics and automation have contributed to progress towards this goal. Specifically, in the area of proteomic identification, it is evident that greater specificity benefits the ability to deal with the high complexity of proteomes. As a result, recent efforts have focused on improvements in separation speed, resolving power and dynamic range, and these methods have generally been based on the combination of separations with mass spectrometry (MS), using correlation of tandem mass spectra with established protein databases or predictions from genome sequence data for identifications.\nAdditionally, modern proteomics research has increasingly taken advantage of the ability of liquid chromatography to identify proteins from their elution time from a chromatographic column. The information gleaned from a liquid chromatograph can be enhanced by identifying the molecule's mass, or mass to charge, by coupling the liquid chromatograph either on line or off line, with a mass spectrometer. Common methods include offline tryptic digestion and subsequent electrophoretic or chromatographic separation with matrix-assisted laser desorption/ionization or electrospray time-of-flight or ion trap mass spectrometry. Capillary electrophoresis, mass spectrometry or liquid chromatography/mass spectrometry coupled online via electrospray interfaces have also been used to analyze tryptic and other digests of complex biological samples such as whole cell lysates and human body fluids. The dynamic range of the mass spectrometer in these methods may be limited when a sample is directly infused by ion suppression in the electrospray and the detector. Further, the dynamic range of Fourier transform ion cyclotron resonance (FTICR) and ion trap mass spectrometers can be limited by the storage capacity within the instrument, although it has been shown that the use of a mass selective quadrupole to selectively load the FTICR cell.\nResearchers attempting to enhance the accuracy of these methods have devised a number of schemes to increase their accuracy. For example, in the paper “Prediction of Chromatographic Retention and Protein Identification in Liquid Chromatography/Mass Spectrometry” Magnus Palmblad, Margareta Ramstrom, Karin E. Markides, Per Hakansson, and Jonas Bergquist, Analytic Chemistry p. 4–9, 2002, the authors describe a method for using the information from liquid separation schemes such as chromatography and electrophoretic methods, to improve peptide mass fingerprinting based on accurate mass measurement. The author's concede that the resolving power and accuracy in chromatographic separations are several orders of magnitude lower than in mass spectrometry, but they contend that the information is complementary in nature and available at negligible computational cost and at no additional experimental cost. Briefly, the method described in the Palmblad paper assigns “retention coefficients” for the 20 amino acids, as well as the number of each amino acid, a term that compensates for void volumes and a delay between sample injection and acquisition of mass spectra. The parameters are then fitted by the least squares method to experimental data from ˜70 BSA peptides of ˜100 HAS and transferrin peptides putatively identified by accurate mass measurement and high relative intensities in the mass spectra. The authors found that “the accuracy of the predictor was found to be 8–10% when “trained” by each of the six BSA and CSF data sets.” While approaches such as that described in the Palmblad paper provide some useful information, their utility is limited by the accuracy of the predictions.\nThus, at the present, there are two major approaches for proteomic analyses. The first one consists of the off-line combination of two-dimensional polyacrylamide electrophoresis (2D-PAGE) with MS. The proteins are first separated in a gel by their pI and mass and then the protein “spots” are enzymatically hydrolysed resulting in peptide mixtures which are analysed by matrix assisted laser desorption ionisation-time of flight (MALDI-TOF) or electrospray (ESI)-MS. Another rapid evolving approach consists of a global proteome-wide enzymatic digestion followed by analysis using on-line 1-D or 2-D liquid chromatography (LC) coupled with ESI-MS. The detection of the peptides is achieved by tandem MS (7,13) or more recently by single stage Fourier transform ion cyclotron resonance (FTICR)-MS, which provides high sensitivity, large dynamic range and high throughput in routine applications by circumventing the need for tandem MS.\nAn aspect of proteomic analysis that has not yet been exploited involves use of the information available from the separations (eg. LC elution time). Indeed, retention time in LC is unique and structurally dependent for a defined experiment (mobile phase composition, stationary phase etc.). If there is a way to predict the LC retention time for a given peptide structure, then this could be used in conjunction with either MS/MS data to improve the confidence of peptide identifications and/or increase the number of peptide identifications, or, with sufficiently high accuracy MS, to reduce the need for MS/MS data (i.e. if the prediction is reliable enough).\nThe idea that chromatographic behaviour of peptides could be predicted based on the amino acid composition is not new. In 1951, Knight and Pardee showed that synthetic peptides retention factor (Rf) values on paper chromatography could be predicted with some accuracy. In 1952, Sanger introduced the problem of isomers by demonstrating that the relationship between Rf and composition was not absolutely accurate since peptides containing the same amino acids but having difference sequences could frequently be separated. More recently, there have been several reports on the prediction of peptide elution times in reversed-phase (RP) or normal phase liquid chromatography. These methods used quantitative structure-chromatographic retention relationships (QSRR's) (e.g. partial least square or multiple linear regression) for the peptide elution time prediction. Casal et al. demonstrated that partial least squares regression provides a better predictive ability with these models using a mixture of 25 small standard peptides. One limitation of these models is that they are most effective for peptides with less than 15–20 amino acid residues.\nAnother approach, based on artificial neural networks (ANNs), has demonstrated better predictive capabilities in several areas of chemistry including: (i) conformational states for small peptides (31), (ii) carbon-13 nuclear magnetic resonance chemical shifts and (iii) the retardation factor or retention time of small molecules in thin layer chromatography, GC and LC. While these advances are significant, until the present invention, those having skill in the art have not yet used ANNs for peptide elution time prediction.\nAccordingly, there remains a need for improved methods for predicting the identity of peptides and proteins."}
{"text": "Analytical chemistry is the analysis of samples to gain an understanding of their chemical composition. The goal of many chemical analysis protocols is to analyze a given sample (e.g., a physiological sample, an environmental sample, a manufacturing sample, etc.) for a variety of different purposes, such as to identify the presence of one or more analytes of interest in the sample, to characterize the makeup of the sample, for example in quality control, etc.\nMany different analytical chemistry protocols have been developed. One broad category of analytical protocols that have been developed is chromatography. Chromatography is a family of analytical chemistry techniques for the separation of mixtures. In chromatography, a sample (the analyte) in a “mobile phase”, often in a stream of solvent, is passed through a “stationary phase”, where the stationary phase is some form of material that will provide resistance between the components of the sample and the material. Usually, each component has a characteristic separation rate that can be used to identify it and thus the composition of the original mixture.\nA chromatograph takes a chemical mixture carried by liquid or gas and separates it into its component parts as a result of differential distributions of the solutes as they flow around or over a stationary liquid or solid phase. Various techniques for the separation of complex mixtures rely on the differential affinities of substances for a gas or liquid mobile medium and for a stationary adsorbing medium through which they pass; such as paper, gelatin, or magnesium silicate gel.\nMany different chromatographic analytical devices have been developed in order to perform various chromatographic protocols. Examples of various chromatographic devices include, but are not limited to: gas chromatography devices, liquid chromatography devices, capillary electrophoresis devices, and supercritical fluid chromatography devices.\nChromatographic devices are typically operated according to an analytical device method, which method is used by a chromatographic device data system (e.g., such as the ChemStation system from Agilent Technologies, Palo Alto, Calif.) to provide all of the setpoints for a device to perform a given sample analysis. As such, an analytical device method generally at least includes instrument control, sample injection and data analysis setpoints. Traditionally, all of the instrument control setpoints for a given method are provided together as a package to a user, e.g., as may be provided in a plurality of selectable complete methods packaged with an analytical device, or as may be imported into the operating data system of a device as a complete method obtained from an outside source. In certain instances, it is possible to import the sample injection and/or data analysis set points as a group into a given data analysis system. In addition, certain chromatographic analytical device data systems provide for editing of one or more parameters of a pre-existing method. However, the inventors are not aware of any product that provides for the ability to selectively import instrument control information into a system that can be used by the system develop a method de novo. Prior solutions have required that the information needed to develop an analysis must be imported in the format defined for that system. For example, current versions of the Agilent ChemStation requires a pre-existing method be imported into the ChemStation methods directory.\nThe access to scientific information has been changed dramatically by the presence of the Internet and by advances in storage media for computers. This improved access has provided electronic access to scientific knowledge in an unprecedented fashion.\nThere is a need in the art to provide for the ability to capitalize on the enhanced access to scientific knowledge in the development of analytical device methods. The present invention satisfies this need."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrochromic display device.\n2. Description of Related Art\nAlong with popularization of an electronic information network, publication in the form of an electronic book, namely electronic publication, has come to be actively performed, in place of publication in the form of a book using a conventional printing technique. As a device which displays electronic information distributed by such a network, for example, a CRT (Cathode Ray Tube) display and a backlight type liquid crystal display have generally been used. However, with respect to a display using these display devices, a place where a user can read the display is more restricted and a weight, size, shape and portability are inferior in terms of handling, compared with a conventional display printed on paper. Further, since these display devices consume large electrical power, when driving them with a battery, a display time is restricted. Moreover, since these display devices are light emission type displays, they may cause too much tiredness when a user fixes the eyes on the devices for a ling time.\nAccordingly, a display device which can solve the problem as described above has been desired. Further, a rewritable display device has been desired. As such display device, one referred to as a paper-like display or electronic paper has been proposed. Specifically, there have hitherto been proposed, for example, a display device of a reflection type liquid crystal system, a display device of an electrophoretic system, a display device of a system in which a dichroic grain is rotated in an electric field, a display device of an electrochromic system (e.g. see Japanese Patent Application Laid-Open Publication No. 2005-84216), and the like.\nHowever, since the electrochromic display device of Japanese Patent Application Laid-Open Publication No. 2005-84216 supplies a scanning signal to a gate wire to sequentially scan pixels so as to rewrite the pixels on a line-by-line basis, in the case of a display device with a large number of lines, such as a large screen and a high definition screen, rewriting its screen takes a long time.\nTherefore, the conventional display device is not suitable for an application requiring a faster rewriting speed and a rapid switching of a screen."}
{"text": "The c-myc proto-oncogene encodes a nuclear phosphoprotein with leucine zipper and helix-loop-helix structural motifs which appears to be important in the molecular biology of normal and abnormal cellular proliferation. Myc is implicated in the control of both differentiation and replication (Cole, Annu. Rev. Genet. 20:361-384 (1986)), and recent reports link myc to apoptotic cell death (Askew et al., Oncogene 6:1915-1922 (1991), Evan et al., Cell, 69:119-125 (1992), and Neiman et al., Proc. Natl. Acad. Sci. USA 88:5857-5861 (1991), each of which is incorporated herein by reference). Myc and its dimerization partner Max form stable heterodimers through their helix-loop-helix and leucine zipper domains and bind specifically to a core \"E box\" CACGTG DNA sequence (Blackwood et al., Science 251:1211-1217 (1991), incorporated herein by reference). Max homodimers may serve as transcriptional repressors, whereas myc/max heterodimers can activate transcription (Kretzner et al., Nature 359:426-429 (1992), incorporated herein by reference). Certain of the biological functions of myc may be mediated by transcriptional regulation of putative target genes.\nDespite recent progress in defining the mechanism of myc action on \"down stream\" events, less progress has been made in defining the proteins regulating the expression of c-myc itself. Both transcriptional and post-transcriptional mechanisms appear to play a role in regulation of c-myc gene expression (Cole, Annu. Rev. Genet. 20:361-384 (1986), Spencer et al., Cancer Res. 56:1-48 (1991), and Marcu et al., Annual Rev. Biochem. 61:809-860 (1992), each of which is incorporated herein by reference). Maintenance of the level of the c-myc MRNA is achieved by regulation of both transcriptional initiation and elongation. Both initiation, and elongation of the c-myc mRNA, depend upon promoter elements which interact specifically with particular nuclear factors (Spencer, Oncogene 5:777-785 (1990) and Spencer et al., Cancer Res. 56:1-48 (1991), each of which is incorporated herein by reference). A general map of mouse and human c-myc transcription elements has been suggested and nuclear factors which bind to these elements have been reported. In certain cases novel cDNA's encoding such factors have been isolated and sequenced including: ZF87 (also called MAZ), a proline-rich six Zn-finger protein binding to ME1a1/ME1a2 elements within P2 promoter of the murine c-myc gene (Pyrc et al., Biochem. 31:4102-4110 (1992) and Bossone et al., Proc. Natl. Acad. Sci. USA, 89:7452-7456 (1992), each of which is incorporated herein by reference); a 37-kDa protein, MBP-1, which appears to be a negative regulator of the human c-myc promoter (Ray et al., Mol. Cell. Biol. 11:2154-2161 (1991), incorporated herein by reference); and nuclease sensitive element protein-1 (NSEP-1) which binds to a region necessary for efficient P2 initiation (Kolluri and Kinniburgh, Nucl. Acids Res. 17:4771 (1991), incorporated herein by reference). In addition, an Rb binding protein E2F which recognizes an E1A-transactivation site in the human c-myc promoter (Thalmeier et al., Genes Dev. 3:527-536 (1989), incorporated herein by reference) has also been cloned (Helin et al., Cell 70:337-350 (1992), incorporated herein by reference).\nThe chicken c-myc 5'-flanking region is at least 10-fold enriched in CpG-pairs compared with total chicken DNA and is presently thought to be a member of the family of CpG-rich islands involved in regulating certain house keeping genes (Bird et al., Nature 321:209-213 (1986), incorporated herein by reference). Overall high GC content (.about.80%) of the 5'-flanking region predicts that most of the potential regulatory DNA elements will be GC-rich. Analysis of DNA-protein interactions within the 5'-flanking region of the chicken c-myc gene revealed multiple GC-rich sequences which specifically interact with nuclear proteins (Lobanenkov et al., Eur. J. Biochem. 159:181-188 (1986), incorporated herein by reference). Proteins binding to one specific region within a hypersensitive site approximately 200 base pairs upstream of the start of transcription have reportedly been analyzed (Lobanenkov et al., Oncogene 5:1743-1753 (1990) and Lobanenkov et al., Gene Reg. and AIDS, Portfolio Publishing Corp., Texas, p. 45-68 (1989), incorporated herein by reference). Three nuclear factors were found that bind to several overlapping sequences within 180-230 bp upstream of the start of transcription. Two of the proteins appear to resemble the transcription factor Sp1, the other is a factor which seems to bind to a GC-rich sequence containing three regularly spaced repeats of the core sequence CCCTC. The CCCTC-binding factor was termed CTCF (Lobanenkov et al., Oncogene 5:1743-1753 (1990) and Lobanenkov et al., Gene Reg. and AIDS, Portfolio Publishing Corp., Texas, p. 45-68 (1989), incorporated herein by reference).\nStudies suggest that during embryonic development the regulatory state of c-myc transcription can determine whether a cell continues to proliferate, or stops, and enters a pathway to terminal differentiation. Failure to properly regulate myc may be one pathway to malignancy. Thus, identifying the suppressor mechanisms by which myc is regulated would provide important reagents and assays useful in the detection of mutants that are indicative of a disease state such as cancer and the development of candidate therapeutic agents can that regulate cell proliferation, for example, inhibiting cell proliferation in cancer on the one hand, or stimulating cell proliferation in a damaged tissue on the other hand. Quite surprisingly, the present invention fulfills these and other related needs."}
{"text": "Generally, transistors of semiconductor devices are classified as NMOS, PMOS, or CMOS according type of channel employed in the transistors. An NMOS type transistor is formed with an N-channel, and a PMOS transistor with a P-channel. In addition, a CMOS (complementary metal oxide silicon) has both NMOS and PMOS, and, thus, both an N-channel and a P-channel are formed therein.\nTo form a CMOS type transistor, an n-well and a p-well are first formed in a horizontal direction on a semiconductor substrate by an ion implantation process, and then shallow trench isolation (STI) is formed. An STI structure prevents a malfunction between neighboring devices by electrically isolating the devices on a semiconductor substrate.\nA well in a semiconductor substrate is classified as a p-well or an n-well according to the type of ions implanted in the well. A p-well is formed on the semiconductor substrate to form an NMOS structure, and an n-well is formed to form a PMOS structure. Subsequently, a gate oxide layer is formed on the semiconductor substrate, and then a polysilicon layer is formed thereon to form a gate stack. The gate stack forms a gate electrode of the NMOS and the PMOS using a photolithography process and an etching process.\nThen, n-type dopants and p-type dopants are respectively implanted into the semiconductor substrate using the gate electrode of the NMOS or the PMOS as an implantation mask. Thus, a source/drain region is formed outward of the gate electrode on an active region of the semiconductor substrate.\nAs described above, the NMOS and the PMOS structures are formed in a horizontal direction in a typical CMOS process and, thus, a CMOS circuit occupies a larger area than does an NMOS circuit or a PMOS circuit. As a result, the CMOS circuit has a drawback in terms of integration.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "The invention relates to pharmaceutically active macrolides, synthesis thereof and intermediates thereto. Halichondrin B is a potent anticancer agent originally isolated from the marine sponge Halichondria okadai, and subsequently found in Axinella sp., Phakellia carteri, and Lisson-dendryx sp. A total synthesis of Halichondrin B was published in 1992 (Aicher, T. D. et al., J. Am. Chem. Soc. 114: 3162-3164). Halichondrin B has demonstrated in vitro inhibition of tubulin polymerization, microtubule assembly, betaS-tubulin crosslinking, GTP and vinblastine binding to tubulin, and tubulin-dependent GTP hydrolysis and has shown in vitro and in vivo anti-cancer properties. Accordingly, there is a need to develop synthetic methods for preparing analogs of Halichondrin B useful as anti-cancer agents."}
{"text": "1. Field of the Invention\nThe present invention relates to a coded modulation system suited for a poor channel in quality impaired by noise, interferences, and/or distortion such as a mobile, portable, or very small aperture terminal (VSAT) satellite communications channel.\n2. Description of the Related Art\nThe duobinary frequency-modulation (DBFM) system with modulation index h=0.5 proposed by P. J. McLane, \"The Viterbi receiver for correlative encoded MSK signals (IEEE Trans. Commun., COM-31, 2, pp. 290-295, Feb. 1983),\" is known as a modulation system that possesses the constant envelope, compact spectrum, and almost the same BER performance as the antipodal modulation systems, i.e., BPSK, QPSK, and MSK.\nThe system combines the duobinary technique with the minimum-shift keying (MSK) and applies a soft-decision Viterbi decoding technique in order to demodulate the signals. However, DBFM has a few disadvantages as follows:\n(1) DBFM necessitates block synchronization between the transmit and receive ends because DBFM applies alternatively two different state-transition tables in order to decode the signals.\n(2) It is difficult for DBFM to employ bit-interleaving because DBFM utilizes the specific behavior of phasor in two-dimensional signal space.\n(3) It is less effective for DBFM to connect tandem in trivial manner such a forward error correction (FEC) codec as a convolutional codec because DBFM already uses up a soft-decision Viterbi decoding of constraint length K=3."}
{"text": "This invention relates to apparatus for preventing the accidental tip-over of free-standing structures such as domestic home appliances.\nIn free-standing electric or gas ranges, the oven door is hinged at its bottom edge and pivots downwardly approximately 90.degree. from a vertical closed position to a horizontal open position. In the open position the door may extend parallel to the floor at a height of several inches to a foot above floor level. In this position it is possible for objects to be placed on the door of sufficient weight to cause the appliance to tilt forward, and possibly to tip completely over.\nOne approach to this problem is described and claimed in commonly assigned U.S. Pat. No. 4,669,695 to Chou. In this arrangement a pair of rigid brackets are mounted to the wall behind the appliance and project forward to be received in corresponding apertures in the rear of the appliance. Each support member is positioned to engage the lower edge of its respective aperture to limit tipping of the appliance when positioned adjacent the wall. The wall engaging portion of each bracket has a length greater than the distance between wall frame members to insure attachment to a frame member. This arrangement has been demonstrated to work satisfactorily. However, the brackets are relatively large to provide the necessary length and rigidity in the portion which engages the appliance and to provide the desired mounting versatility, rendering the approach relatively costly.\nU.S. Pat. No. 4,754,948 to Casciani discloses an alternative approach in which a pair of U-shaped brackets are disposed proximate to the intersection of wall and floor with the vertically extending bight attached to the wall and one leg attached to the floor. The free leg projects forward to extend into an aperture formed in the rear cabinet wall of the appliance. The U-shaped brackets must be rigid enough and long enough to prevent the tipping movement of the appliance from causing it to slide off of the retaining arm. In addition, a pair of brackets is required. In an alternative arrangement therein disclosed, a pair of rearwardly extending projections are secured to the rear of the appliance to extend into corresponding holes drilled in the wall behind the appliance. This use of projections extending from the appliance presents obvious alignment problems rendering installation difficult.\nIn view of the limitations of the foregoing prior art, it is a primary object of the present invention to provide an improved anti-tip apparatus for appliances which is easy to install initially and which facilitates removal and reinstallation such as for cleaning around or servicing the appliance, and which is relatively inexpensive."}
{"text": "Recent years have seen development of fuel cells that have good power generation efficiency and do not emit carbon dioxide from the viewpoint of preserving the global environment. This fuel cell generates power by causing hydrogen and oxygen to react with each other. A basic structure of a fuel cell resembles a sandwich and is constituted by an electrolyte membrane (i.e., ion exchange membrane), two electrodes (i.e., a fuel electrode and an air electrode), a diffusion layer for diffusing hydrogen and oxygen (air), and two separators. Phosphoric-acid fuel cells, molten carbonate fuel cells, solid-oxide fuel cells, alkaline fuel cells, proton-exchange membrane fuel cells, and the like have been developed in accordance with the type of electrolyte used.\nOf these fuel cells, proton-exchange membrane fuel cells in particular have the following advantages over molten carbonate fuel cells, phosphoric-acid fuel cells, and the like:\n(a) Operating temperature is significantly low, i.e., about 80° C.\n(b) Weight- and size-reduction of the fuel cell main body is possible.\n(c) The time taken for start-up is short and fuel efficiency and output density are high.\nAccordingly, proton-exchange membrane fuel cells are one of the most prospective fuel cells today, for onboard power supplies for electric vehicles and portable and compact dispersed power systems for household use (stationary type compact electric generator).\nA proton-exchange membrane fuel cell is based on the principle of extracting power from hydrogen and oxygen through a polymer membrane and has a structure shown in FIG. 1, in which a membrane-electrode assembly 1 is sandwiched by gas diffusion layers 2 and 3 such as carbon cloths and these form a single constitutional element (also known as a single cell). Electromotive force is generated between the separators 4 and 5.\nThe membrane-electrode assembly 1 is also known as MEA (Membrane-Electrode Assembly) and is made by integrating a polymer membrane and an electrode material such as carbon black carrying a platinum catalyst, the electrode material being provided on front and back surfaces of the polymer membrane. The thickness of the membrane-electrode assembly 1 is several ten to several hundred micrometers. The gas diffusion layers 2 and 3 are frequently integrated with the membrane-electrode assembly 1.\nWhen proton-exchange membrane fuel cells are applied to the usages described above, several ten to several hundred single cells described above are connected in series to form a fuel cell stack, and the fuel cell stack is used.\nThe separators 4 and 5 are required to have\n(A) a function of a separator that separates between single cells each other, as well as\n(B) a function of an electric conductor that carries electrons generated;\n(C) a function of a channel for oxygen (air) and hydrogen (air channels 6 and hydrogen channels 7 in FIG. 1); and\n(D) a function of a discharge channel for discharging water and gas generated (air channels 6 and hydrogen channels 7 also serve as this discharge channel).\nIn order to use a proton-exchange membrane fuel cell in practical application, separators having good durability and conductivity must be used.\nThe durability expected is about 5000 hours for fuel cells for electric vehicles and about 40000 hours for stationary type electric generators used as compact dispersed power systems for household use and the like.\nProton-exchange membrane fuel cells that have been put to practice hitherto use carbon materials as separators. However, since the separators using carbon materials are susceptible to fracture upon impact, they have the drawbacks that not only the size-reduction is difficult but also the process cost for forming channels is high. In particular, the cost problem has been the largest impediment for spread of fuel cells.\nIn response, attempts have been made to use a metal material, in particular, stainless steel, instead of carbon materials as the material for separators.\nThe operating environment the separators are exposed to are characteristic in that the environment is acidic and has a high temperature of 70° C. or higher and the expected potential range is as wide as from about 0 V vs SHE to 1.0 V vs SHE or higher (hereinafter all potentials are versus SHE and simply denoted as V). In order to use stainless steel, the corrosion resistance in the expected potential range needs to be improved. In particular, at and near 1.0 V, transpassive dissolution of Cr, which is the main element of the stainless steel, occurs and thus it is difficult to maintain corrosion resistance solely by Cr on one hand. On the other hand, Cr is primarily responsible for maintaining the corrosion resistance at 0.6 V or less. Thus, according to the conventional art, the corrosion resistance could not be maintained in a wide potential range from a low potential to a high potential.\nFor example, patent document 1 discloses a stainless steel for a separator in which the corrosion resistance is improved from the composition aspect by increasing the Cr and Mo contents.\nPatent document 2 discloses a method for producing a separator for a low-temperature-type fuel cell characterized in that a stainless steel sheet containing 0.5 mass % or more of Cu is subjected to alternation electrolytic etching of alternately performing anodic electrolyzation at a potential of +0.5 V or more and cathodic electrolyzation at a potential between −0.2 V and −0.8 V in an aqueous solution of ferric chloride.\nPatent document 3 discloses a stainless steel conductive part and method of producing the same that has excellent conductivity and low contact electrical resistance formed by modifying a passive film on a stainless steel surface by injecting fluorine in the passive film."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of preparing a catalyst for polymerization of aliphatic polycarbonates and a method of polymerizing an aliphatic polycarbonate, and more particularly, to a method of preparing a catalyst for polymerization of aliphatic polycarbonates exhibiting high catalyst activity using an amphiphilic block copolymer.\n2. Background of the Invention\nCarbon dioxide from industrial activities, among atmospheric pollutants, has been known as one reason for climatic change according to the UNFCCC, so various studies to reduce the amount of carbon dioxide produced have been undertaken all around the world. Therefore, in order to protect the environment and to use carbon dioxide, a method in which an epoxide reacts with carbon dioxide as a carbon source in the presence of a zinc-included catalyst to prepare an aliphatic polycarbonate has attracted attention.\nThe aliphatic polycarbonate is able to form a film or a particle, and has uses in many areas such as for ceramic binders, evaporative pattern casting, and adhesives. However, this method has a low yield because of low carbon dioxide reactivity. Accordingly, it is difficult to use industrially, so it is required to prepare a catalyst exhibiting high efficiency for increasing the yield of the aliphatic polycarbonate.\nInoue teaches a method of polycarbonate production from carbon dioxide and epoxide in U.S. Pat. No. 3,585,168.\nThe Inoue catalyst system was prepared by the reaction of a diethylzinc catalyst with materials containing active hydrogen compounds, e.g., water, dicarboxylic acid, or dihydric phenols, and the typical catalyst productivities ranged from 2.0 to 10.0 grams of polymer per gram of catalyst used. The catalyst has shortcomings associated with use and storage, because of stability and sensitivity to moisture and to other catalyst poisons, and it has a low yield, so it has been required to study other catalyst systems.\nZinc dicarboxylic acid esters (Polymer J. 13(4), 407(1981)) reported by Soga have also been described as effective catalysts for copolymerization of carbon dioxide and propylene oxide, and since these are stable materials with none of the handling problems associated with diethylzinc, they represent interesting candidates for a practical commercial catalyst system.\nMotika (U.S. Pat. No. 5,026,676) teaches a method for preparing zinc dicarboxylic acid ester in which zinc oxide reacts with dicarboxylic acid in the presence of an organic solvent. Glutaric acid and adipic acid produced catalysts with higher activity than the known zinc dicarboxylic acid ester catalysts, and the catalyst production is about 2 to 26 grams of the aliphatic polycarbonate per gram of catalyst.\nZinc dicarboxylic acid ester is a heterogeneous catalyst of which catalyst activity depends on its surface structure, so various processes to form a surface structure that gives high catalyst activity have been proposed. For example, Kawachi (U.S. Pat. No. 4,981,948) teaches a method of producing zinc dicarboxylic acid ester with a mechanical pulverization treatment such as with a ball mill, and the yield of the aliphatic polycarbonate is 8.1 to 34.2 gram per gram of the catalyst.\nIn addition, a method of preparing zinc dicarboxylic acid ester as a homogeneous catalyst is disclosed in U.S. Pat. No. 4,783,445, by Sun. This method includes a reaction of zinc salts with a dicarboxylic acid monoester in an organic solvent, and the yield is 5.1 to 12.4 gram per gram of the catalyst."}
{"text": "A strut may be connected to a vehicle using a top mount."}
{"text": "Manufacturing a shoe typically requires a number of assembly steps, such as cutting, forming, assembling, adhering, and/or stitching several shoe parts together. Some methods of completing these steps, such as those that rely heavily on manual execution, may be resource intensive and may have a high rate of variability."}
{"text": "The present invention relates to a dynamic type RAM (random access memory) and, more particularly, to a technique which may be effectively utilized for a semiconductor integrated circuit device having a dynamic type RAM formed using one-element type dynamic memory cells each consisting of a data storing capacitor and an address selecting MOSFET.\nAs a result of the increase in the storage capacity of dynamic type RAMs, the chip size of the semiconductor substrate has become larger, and the distributed resistance of word lines which are formed from a polycrystalline silicon or the like has become a serious problem. As one of the solutions to this problem, the Al shunt method is known in which word lines are divided into appropriate lengths to define divided word lines which are coupled to word lines which are formed from an aluminum layer having a relatively large conductivity. The Al shunt method is described, for example, in \"Digest of Technical Papers\" of ISSCC, SESSION-XVI, February 1983, pp. 226-227."}
{"text": "1. Field of the Invention\nThis invention relates to a misfire-detecting system for internal combustion engines, and more particularly to a misfire-detecting system of this kind, which is adapted to detect a misfire attributable to the fuel supply system.\n2. Prior Art\nIn an internal combustion engine having spark plugs, a misfire can occur, in which normal ignition does not take place at one or more of the spark plugs. Misfires are largely classified into ones attributable to the fuel supply system and ones attributable to the ignition system. Misfires attributable to the fuel supply system are caused by the supply of a lean mixture or a rich mixture to the engine, while misfires attributable to the ignition system are caused by failure to spark (so-called mis-sparking), i.e. normal spark discharge does not take place at the spark plug, for example, due to smoking or wetting of the spark plug with fuel, particularly adhesion of carbon in the fuel or unburnt fuel to the spark plug, or abnormality in the sparking voltage supply system.\nThe present assignee has already proposed a misfire-detecting system for detecting misfires attributable to the fuel supply system, which comprises sparking voltage detecting means which detects sparking voltage, i.e. voltage across electrodes of the spark plug, and misfire-determining means which determines that a misfire has occurred when a time period over which the detected value of the sparking voltage exceeds a predetermined reference value (Japanese Patent Application No. 3(1991)-326507 and corresponding U.S. Ser. No. 07/846,238 filed Mar. 5, 1992).\nIn the above proposed system, the time period over which the detected value of sparking voltage exceeds the predetermined reference value corresponds to a time period over which a predetermined amount of electric charge or more is stored in floating capacitance in the vicinity of the spark plug. Depending upon the behavior of discharge caused by sparking of the spark plug, the charge can be discharged within a short time period even if a misfire has occurred. This phenomenon can take place when the sparking voltage assumes a considerably high voltage value at the end of an inductive discharge caused by sparking, due to occurrence of a misfire. In such an event, discharge again takes place between the electrodes of the spark plug and terminates within a short time period so that the misfire is not detected to have occurred.\nFurther, even in the case where the spark plug has just started smoking due to adhesion of carbon, etc. to the electrodes of the spark plug and hence has decreased insulation resistance between the electrodes, or in the case where the spark plug has just started recovering from its smoking state due to its own purifying action, dielectric breakdown is likely to occur so that there is no significant difference in the time period over which the sparking voltage exceeds the predetermined reference value between when normal firing has occurred and when a misfire has occurred."}
{"text": "In many industries, such as health care, banking, and finance, the protection and ensuring of a person's privacy is often considered a top priority. In certain cases, ensuring privacy may also be mandated by law. Thus, various health care services, providers, and stakeholders (i.e. a doctor, nurse, or the like) may implement processes and methods that ensure that a person's Protected Health Information (PHI) is secure and shared with the actual person whose health the information is about. In the banking, credit, and finance industries, both Protected Individual Information (PII) and Protected Credit Information (PCI) are secured and shared with the person to whom the information belongs. PHI, PII, and PCI are of great value and can be targets for theft by hackers who can resell the information.\nIn the present health care industry, various health care services exist that allow a care giver, a hospital, or other relevant health care provider stakeholder to provide services to a patient. In providing health care, the various stakeholders may access private information about a patient or a group of patients and their Protected Health Information (PHI) and Protected Identity Information (PII) in order to provide the health care services. Other service-oriented industries are similar, such as a teller accessing the PII or Protected Credit Information (PCI) of a bank customer, or a credit authorization agent verifying the PII and PCI of a credit card user to permit a financial transaction to be completed.\nThe services to retrieve and access PHI, PCI and PII may be guarded by one or more security firewalls, each of which requires an authentication process to get beyond to access the information. The authentication process may be referred to in healthcare or other industries as authentication, identity authentication, or identity verification."}
{"text": "This invention relates to medical diagnostic x-ray apparatus, particularly apparatus for performing angiography.\nIn angiography procedures it is frequently necessary to obtain simultaneous x-ray views of the blood vessels in two different directions, such as in the postero-anterior direction and in the lateral direction. Apparatus which permits postero-anterior views is shown in U.S. Ser. No. 202,094, filed by Stivender et al. on Oct. 31, 1980 now U.S. Pat. No. 4,358,856 and assigned to the owner of the present invention. That application is incorporated herein by reference in its entirety to explain the construction and operation of the L-U arm apparatus for taking such views.\nTwo approaches are currently used for taking simultaneous lateral views. In the first approach, an x-ray source hung from the ceiling is positioned on one side of the patient and a freestanding x-ray detection device is positioned on the other side of the patient. As is well known, the freestanding detector and its associated electrical cables prevent the physician from moving freely around the patient and can also interfere with the source or detector for taking postero-anterior views. Another deficiency of this apparatus is that the source and detector for lateral views must be aligned manually.\nA second approach is shown in FIGS. 12 and 13, illustrating two prior art devices. In the device shown in FIG. 12, a single structural member 20 carries an x-ray source 22 and an electronic image intensifier 24 at its respective ends. Member 20 is supported by a brace 26 pivotally mounted to an overhead support 28 for rotating the pattern of radiation passing from source 22 to detector 24 about a vertical axis without disturbing the relative alignment of source and detector. In the embodiment shown in FIG. 13, source 22 and detector 24 are rigidly mounted to the respective ends of a rigid C-shaped member 30 received in a guide 32, which again is pivotally mounted to an overhead support 34. In this embodiment, member 30 can be rotated as before, or can be driven in either direction through guide 32 to rotate the pattern of x-rays passing between source 22 and detector 24 about the longitudinal axis of a patient.\nThe devices of FIGS. 12 and 13 seriously interfere with access to the patient by the physician, and when in motion can present a hazard to the patient and those working around the patient. Also, since in both prior art embodiments the mass of the x-ray tube and image intensifier is supported at a single point between them, support members 20 and 30 are prone to gravitational and inertial bending moments and oscillations which complicate the problem of aiming source 22 at detector 24. Furthermore, such devices can be disturbing to the patient, who is encircled by machinery. These structures also are difficult or impossible to move out of the way when they are not in use, as the entire assembly must be moved as a unit and cannot be retracted or collapsed to provide head room."}
{"text": "An investigation has been conducted on use of, for example, an alloy as a hydrogen-storing material expected to be applicable to a hydrogen-storing system in a fuel-cell vehicle or the like. However, when a hydrogen-storing alloy is used, its hydrogen storage capacity is insufficient. In addition, depending on kinds of metal to be used, the alloy not only has low durability but also involves problems in terms of price and reserve.\nOn the other hand, use of a carbon material that raises no concern about exhaustion of resources and is relatively inexpensive has been investigated. For example, Patent Literature 1 describes a carbon material whose hydrogen storage quantity is increased by expanding an average distance between carbon layers to 0.5 nm or more to cause the carbon layers to hold hydrogen therebetween. In addition, Patent Literature 2 describes an activated carbon material with an increased hydrogen storage quantity as a result of possession of a pore diameter of 0.3 nm or more and 1.5 nm or less."}
{"text": "Gypsum-based building products are commonly used in construction. Wallboard made of gypsum is fire retardant and can be used in the construction of walls of almost any shape. It is used primarily as an interior wall and ceiling product. Gypsum has sound-deadening properties. It is relatively easily patched or replaced if it becomes damaged. There are a variety of decorative finishes that can be applied to the wallboard, including paint and wallpaper. Even with all of these advantages, it is still a relatively inexpensive building material.\nGypsum is also known as calcium sulfate dihydrate, terra alba or landplaster. Plaster of Paris is also known as calcined gypsum, stucco, calcium sulfate semihydrate, calcium sulfate half-hydrate or calcium sulfate hemihydrate. Synthetic gypsum, for example, that which is a byproduct of flue gas desulfurization processes from power plants, may also be used. When it is mined, raw gypsum is generally found in the dihydrate form. In this form, there are two water molecules associated with each molecule of calcium sulfate. To produce the hemihydrate form, the gypsum is calcined to drive off some of the water of hydration by the following equation:CaSO4.2H2O→CaSO4.1/2H2O+3/2H2O\nA number of useful gypsum products can be made by mixing the stucco with water and permitting it to set by allowing the calcium sulfate hemihydrate to react with water to convert the hemihydrate into a matrix of interlocking calcium sulfate dihydrate crystals. As the matrix forms, the product slurry becomes firm and holds a desired shape. Excess water must then be removed from the product by drying.\nSignificant amounts of energy are expended in the process of making gypsum articles. Landplaster is calcined to make stucco by heating it to drive off water of hydration. Later the water is replaced as the gypsum sets by hydration of the hemihydrate to the dihydrate form. Excess water used to fluidize the slurry is then driven from the set article by drying it in an oven or a kiln. Thus, reducing the amount of water needed to fluidize the slurry turns into a monetary savings when fuel requirements are decreased. Additional fuel savings would result if the amount of material that required calcining were reduced.\nAttempts have been made to reduce the amount of water used to make a fluid slurry using dispersants. Polycarboxylate superplasticizers are very effective in allowing water reduction and where water reduction results in increased density, a strength increase is achieved. These materials are relatively expensive. When used in large doses, polycarboxylate dispersants can be one of the single, most expensive additives in making gypsum products. The high price of this component can overcome the narrow margins afforded these products in a highly competitive marketplace.\nAnother disadvantage associated with polycarboxylate dispersants is the retardation of the setting reaction. Gypsum board is made on high-speed production lines where the slurry is mixed, poured, shaped and dried in a matter of minutes. The board must be able to hold its shape to be moved from one conveyor line to another to put the board into the kiln. Damage can occur if the boards have not attained a minimum green strength by the time they are cut to length and handled during the manufacturing process. If the board line has to be slowed down because the board is not sufficiently set to move on to the next step in the process, production costs are driven up, resulting in an economically uncompetitive product.\nModifiers have been found that increase the efficacy of the dispersant in fluidizing the slurry, allowing the modifier to replace a portion of the expensive dispersant while still reducing water demand. However, it has been found that the modifier does not work consistently, depending on how and when it is added to the slurry. Thus, there is a need for a delivery vehicle to carry the modifier to the slurry in a manner that allows it to perform consistently so that the amount of dispersant can be reduced.\nThe use of fillers that are easily fluidizable in water have been considered as another method of reducing fuel demand. However, one of the important properties of gypsum products, and especially gypsum panels or wallboard, is its fire resistance. Calcium sulfate dihydrate is approximately 20% water by weight. Replacing a portion of the calcined gypsum with fillers that are less fire retardant diminishes this property in the finished product. Many fillers also reduce the compressive strength and the nail pull strength of wallboard.\nLandplaster has been used as a filler in gypsum products. It is also fire retardant, inexpensive, readily available and reduces the amount of calcined gypsum that is needed, but it also has disadvantages. Calcium sulfate dihydrate used in sufficient quantities to act as a filler also acts as a set accelerator for the hemihydrate by providing seed crystals that start the crystallization process more quickly. This leads to premature stiffening of the slurry.\nThus there is a need in the art for a filler for use in gypsum articles, particularly wallboard, that reduces fuel consumption by replacing calcined gypsum, by reducing the amount of water driven from the set product or both. The filler should have fire retardancy approximately equal to set gypsum and it should be inexpensive, readily available and should not decrease the strength of the finished product.\nThus, there is a need in the art to reduce the dosage of dispersants used in a gypsum slurry while maintaining flowability of the slurry. Reduction in dispersant use would result in saving of costs spent on the dispersant and would reduce adverse side effects, such as set retardation."}
{"text": "Objective lens assemblies are commonly used in microscopes, telescopes, cameras and other devices for gathering light from an object being observed and focusing the light to form an image of the object. Objective lens assemblies that operate in visible spectrum of light are quite common.\nCurrently, the applicant of the present invention is developing a microscope that operates in the mid infrared (“MIR”) light spectrum. Unfortunately, existing objective lens assemblies do not provide sufficient performance in the MIR light spectrum."}
{"text": "The present disclosure relates to an analysis system for determining an analyte in a body fluid having a disposable integrated sample acquisition and analysis element and having a reusable analysis instrument.\nFor diagnostic, purposes, small quantities of body fluids, such as blood, are taken from a body part. For this purpose, piercing devices having lancets are typically used, such devices generate a wound in the body part, for example, in the finger or in the ear lobe. The piercing systems are implemented in such a manner that they may also be used by laymen.\nHowever, when determining an analyte in the body fluid, a procedure in multiple steps is required. Firstly, a wound must be generated in the body part using a piercing system, which comprises a piercing device and a lancet. A body fluid, such as, for example, blood, then exits from the wound. In a further step, the fluid must be received by a test element and supplied to an analysis system, which determines the desired analyte in the body fluid.\nThis procedure is complex, in particular for diabetics, who must determine the glucose content in the blood multiple times a day. Therefore, in addition to pain-free piercing as much as possible, increased operating comfort is also required.\nIn order to come one step closer to the desired operating comfort, integrated analysis systems have been proposed in the prior art, which also comprise an analysis unit in addition to a piercing device. An analysis system is known from WO 2006/027101 A1, in which, after the generation of the wound in the body part, the piercing device is moved away from an opening of the analysis system and an analysis unit is moved to the opening of the system in a second step, so that blood exiting from the wound can be received by the analysis unit.\nIn order to improve the handling ability and the comfort of the analysis systems, analysis systems have been developed which propose test elements and/or test sensors having an integrated lancet.\nFor example, an analysis system having a test sensor with an integrated lancet, in which an analyte in a body fluid is determined by an electrochemical measurement is known from WO2006/092281. The disposable test sensor comprises a test strip, on whose top side a capillary channel is provided, which is used to transfer the received body fluid to test electrodes, which are also located on the top side of the test strip, in order to determine an analyte in the body fluid. A lancet is located on the bottom side of the test strip, which is movably mounted relative to the test strip. The lancet is enclosed by a sterile envelope, from which it exits before piercing into the body part.\nFor the sample acquisition, the test sensor is guided into the vicinity of the opening of the analysis system. The lancet is then moved forward in the puncture direction until it exits from the opening of the receptacle system and generates a wound in a body part which is pressed against the opening. After the puncture, the lancet is retracted again until it is positioned in its sterile protective envelope again. The blood exiting from the puncture wound and/or the exiting body fluid is suctioned in by the capillary channel on the top side of the test strip and finally reaches the electrodes, so that an analyte in the body fluid may be determined electrochemically.\nA system of this type has the disadvantage that two separate handling steps are necessary in order to, on the one hand, ensure the puncture in the skin and, on the other hand, ensure the transfer of a sufficiently large blood sample onto the test element.\nThis may be performed manually, after execution of the puncture, the analysis system being removed from the body part and the body part subsequently being “milked” in order to promote the escape of blood from the wound. As soon as a sufficiently large quantity of blood has exited from the wound, the analysis system is guided back to the wound manually in order to suction the sample into the capillary channel. This requires cumbersome handling by the patient and is difficult in particular for older people, who are frequently affected by diabetes. Alternatively to the manual procedure and the operation comprising multiple steps, the “milking” may be mechanized by the analysis system itself. For this purpose, systems having a so-called finger cone have been proposed. After the piercing procedure, the expression of a liquid sample is caused by pressure on the cone. For this purpose, the manual handling is thus replaced by a corresponding instrument function, which requires significant design effort and makes the instruments more costly, however.\nAs a third, rather theoretical possibility, which is also opposed by the demand for piercing with as little pain as possible, the lancet may be pierced so deeply into the body part that a sufficiently large blood droplet exits without additional measures such as manual or mechanical milking. The pain connected with the deep piercing is so great, however, that a system of this type is unsuitable for practice.\nTherefore, there is a need for an analysis system which is distinguishable from the prior art by a high comfort for the user, in particular in that he only has to hold the analysis system on the body part once in order to determine the desired analyte. In addition, the analysis system will be easy to operate and will be based on a simple and cost-effective design."}
{"text": "During wire-bonding of semiconductor devices, wherein electrical connections are made between bond pads of dice and/or substrates on which they are attached, it is common to utilize a wire clamp to feed a roll of bonding wire towards a bonding site. The clamp is opened to allow wire to feed through during threading of the wire through a capillary and thereafter closed to control the wire. The wire clamp may also be used to hold the wire in position during the making of a first bond and a second bond on the die and/or substrate. The clamp is further commonly used to enable looping of a length of bonding wire between electrical contact points on the die and/or substrate, and/or to pull wires from bonds after the bonds have been made.\nThe clamp typically comprises a movable arm or member, and a fixed arm or member. The movable arm is opened and closed by a solenoid or a linear motor, and is usually urged towards the fixed arm by a spring or the motor. The bonding wire is very fine, to the order of 1 mil or less. Thus the wire is easily broken if subjected to excessive force. It is important that a clamping force exerted by the wire clamp is sufficient to grasp the wire, but not too high so as to cause abnormal deformation or to break the wire. It is also important that a gap between the movable and fixed arms is sufficient for the wire to pass through, and yet not be so large as compared to the size of the wire when opened so that the clamping force cannot be easily or reliably controlled.\nIn view of the above, it is usually necessary to calibrate a wire clamp prior to using it. Prior art devices for calibrating wire clamps have been devised, but such prior art devices involve too much human intervention and have become less effective as the diameters of bonding wires decrease together with decreasing dimensions of semiconductor packages.\nFor example, in order to measure clamping force, a gram gauge (see FIG. 1) has been used in the prior art. The gram gauge has a deflectable lever, which measures a deflecting force exerted on the lever by a movable wire clamp member. The gram gauge functions in much the same way as a conventional weighing scale. However, a spring which biases the lever during deflection is not sufficiently sensitive where the clamping force is small and only deflects the deflectable lever minimally.\nIn order to measure a gap between clamping members, a thin gauge sheet of a known thickness may be inserted between the clamping members (see FIG. 2). If the gauge sheet cannot fit into the gap, it means that the distance between the clamp members is smaller than the thickness of the gauge sheet. If the gauge sheet can fit into the gap with space to spare, then the distance is much larger than the thickness of the gauge sheet. The distance between the clamp members should be adjusted so that the gauge sheet just fits into the gap. Thus, this method is based on trial-and-error, and the error margin gets larger as the distance between clamp members (for smaller diameters of wires) gets smaller. In the event, this method does not offer sufficient accuracy for thinner wires for smaller semiconductor packages."}
{"text": "1. Field of the Invention\nThe present invention relates to a circuit breaker used in a power semiconductor device, such as an inverter that passes large current, and specifically to a small and low-cost circuit breaker that can reliably break a circuit regardless of operating conditions, and has a small wiring loss.\n2. Background Art\nThere is a case wherein an element of an inverter for driving a motor installed in a hybrid motor vehicle or the like is broken to be a short-circuited state, and regenerative current flows back from a motor rotated by an engine power. In order to prevent the flow of a large current into the circuit in such abnormality, a circuit breaker has been used. As a circuit breaker, a device wherein a fuse is broken by the heat of a heater has been proposed (see, for example, Japanese Patent Laid-Open No. 6-119858).\nIn a conventional circuit breaker, a heater was serially inserted in a current path, and a fuse was broken when a larger current than in normal operation flows in abnormality. However, in the case of a hybrid vehicle or the like, since there was no significant difference in the current value for driving a motor in normal operation and the regenerative current value from the motor in abnormality, the fuse cannot be reliably broken. In addition, since an electric resistor serially inserted in a current path was used as the heater, there was a problem of large wiring loss. If the fuse was substituted by a relay, although the circuit was reliably broken, there were problems of contact resistance, costs, and a space."}
{"text": "(1) Field of the Invention\nThe present invention relates generally to the field of sonar signal processing and, more particularly, to determining whether d-dimensional data sets are random or non-random in nature.\n(2) Description of the Prior Art\nNaval sonar systems require that signals be classified according to structure; i.e., periodic, transient, random or chaotic. For instance, in many cases it may be highly desirable and/or critical to know whether data received by a sonar system is simply random noise, which may be a false alarm, or is more likely due to detection of sound energy emitted from a submarine or other vessel of interest. In the study of nonlinear dynamics analysis, scientists, in a search for “chaos” in signals or other physical measurements, often resort to “embedding dimensions analysis,” or “phase-space portrait analysis.” One method of finding chaos is by selecting the appropriate time-delay close to the first “zero-crossing” of the autocorrelation function, and then performing delay plot analyses. Other methods for detection of spatial randomness are based on an approach sometimes known as “box counting” and/or “box counting enumerative” models. Other methods such as power spectral density (PSD) techniques may be employed in naval sonar systems. Methods such as these may be discussed in the subsequently listed patents and/or the above-cited related patent applications which are hereby incorporated by reference and may also be discussed in patents and/or applications by the inventors of the above-cited related patent applications and/or subsequently listed patents.\nIt is also noted that recent research has revealed a critical need for highly sparse data set time distribution analysis methods and apparatus separate and apart from those adapted for treating large sample distributions. It is well known that large sample methods often fail when applied to small sample distributions, but that the same is not necessarily true for small sample methods applied to large data sets. Very small data set distributions may be defined as those with less than about ten (10) to thirty (30) measurement (data) points.\nExamples of exemplary patents related to the general field of the endeavor of analysis of sonar signals include:\nU.S. Pat. No. 5,675,553, issued Oct. 7, 1997, to O'Brien, Jr. et al., discloses a method for filling in missing data intelligence in a quantified time-dependent data signal that is generated by, e.g., an underwater acoustic sensing device. In accordance with one embodiment of the invention, this quantified time-dependent data signal is analyzed to determine the number and location of any intervals of missing data, i.e., gaps in the time series data signal caused by noise in the sensing equipment or the local environment. The quantified time-dependent data signal is also modified by a low pass filter to remove any undesirable high frequency noise components within the signal. A plurality of mathematical models are then individually tested to derive an optimum regression curve for that model, relative to a selected portion of the signal data immediately preceding each previously identified data gap. The aforesaid selected portion is empirically determined on the basis of a data base of signal values compiled from actual undersea propagated signals received in cases of known target motion scenarios. An optimum regression curve is that regression curve, linear or nonlinear, for which a mathematical convergence of the model is achieved. Convergence of the model is determined by application of a smallest root-mean-square analysis to each of the plurality of models tested. Once a model possessing the smallest root-mean-square value is derived from among the plurality of models tested, that optimum model is then selected, recorded, and stored for use in filling the data gap. This process is then repeated for each subsequent data gap until all of the identified data gaps are filled.\nU.S. Pat. No. 5,703,906, issued Dec. 30, 1997, to O'Brien, Jr. et al., discloses a signal processing system which processes a digital signal, generally in response to an analog signal which includes a noise component and possibly also an information component representing three mutually orthogonal items of measurement information represented as a sample point in a symbolic Cartesian three-dimensional spatial reference system. A noise likelihood determination sub-system receives the digital signal and generates a random noise assessment of whether or not the digital signal comprises solely random noise, and if not, generates an assessment of degree-of-randomness. The noise likelihood determination system controls the operation of an information processing sub-system for extracting the information component in response to the random noise assessment or a combination of the random noise assessment and the degree-of-randomness assessment. The information processing system is illustrated as combat control equipment for submarine warfare, which utilizes a sonar signal produced by a towed linear transducer array, and whose mode operation employs three orthogonally related dimensions of data, namely: (i) clock time associated with the interval of time over which the sample point measurements are taken, (ii) conical angle representing bearing of a passive sonar contact derived from the signal produced by the towed array, and (iii) a frequency characteristic of the sonar signal.\nU.S. Pat. No. 5,966,414, issued Oct. 12, 1999, to Francis J. O'Brien, Jr., discloses a signal processing system which processes a digital signal generated in response to an analog signal which includes a noise component and possibly also an information component. An information processing sub-system receives said digital signal and processes it to extract the information component. A noise likelihood determination sub-system receives the digital signal and generates a random noise assessment that the digital signal comprises solely random noise, and controls the operation of the information processing sub-system in response to the random noise assessment.\nU.S. Pat. No. 5,781,460, issued Jul. 14, 1998, to Nguyen et al., discloses a chaotic signal processing system which receives an input signal from a sensor in a chaotic environment and performs a processing operation in connection therewith to provide an output useful in identifying one of a plurality of chaotic processes in the chaotic environment. The chaotic signal processing system comprises an input section, a processing section and a control section. The input section is responsive to input data selection information for providing a digital data stream selectively representative of the input signal provided by the sensor or a synthetic input representative of a selected chaotic process. The processing section includes a plurality of processing modules each for receiving the digital data stream from the input means and for generating therefrom an output useful in identifying one of a plurality of chaotic processes. The processing section is responsive to processing selection information to select one of the plurality of processing modules to provide the output. The control module generates the input data selection information and the processing selection information in response to inputs provided by an operator.\nU.S. Pat. No. 5,963,591, issued Oct. 5, 1999, to O'Brien, Jr. et al., discloses a signal processing system which processes a digital signal generally in response to an analog signal which includes a noise component and possibly also an information component representing four mutually orthogonal items of measurement information representable as a sample point in a symbolic four-dimensional hyperspatial reference system. An information processing and decision sub-system receives said digital signal and processes it to extract the information component. A noise likelihood determination sub-system receives the digital signal and generates a random noise assessment of whether or not the digital signal comprises solely random noise, and if not, generates an assessment of degree-of-randomness. The noise likelihood determination system controls whether or not the information processing and decision sub-system is used, in response to one or both of these generated outputs. One prospective practical application of the invention is the performance of a triage function upon signals from sonar receivers aboard naval submarines, to determine suitability of the signal for feeding to a subsequent contact localization and motion analysis (CLMA) stage.\nU.S. Pat. No. 6,397,234, issued May 28, 2002, to O'Brien, Jr. et al., discloses a method and apparatus are provided for automatically characterizing the spatial arrangement among the data points of a time series distribution in a data processing system wherein the classification of said time series distribution is required. The method and apparatus utilize a grid in Cartesian coordinates to determine (1) the number of cells in the grid containing at least-one input data point of the time series distribution; (2) the expected number of cells which would contain at least one data point in a random distribution in said grid; and (3) an upper and lower probability of false alarm above and below said expected value utilizing a discrete binomial probability relationship in order to analyze the randomness characteristic of the input time series distribution. A labeling device also is provided to label the time series distribution as either random or nonrandom.\nU.S. Pat. No. 5,144,595, issued Sep. 1, 1992, to Graham et al., discloses an adaptive statistical filter providing improved performance target motion analysis noise discrimination includes a bank of parallel Kalman filters. Each filter estimates a statistic vector of specific order, which in the exemplary third order bank of filters of the preferred embodiment, respectively constitute coefficients of a constant, linear and quadratic fit. In addition, each filter provides a sum-of-squares residuals performance index. A sequential comparator is disclosed that performs a likelihood ratio test performed pairwise for a given model order and the next lowest, which indicates whether the tested model orders provide significant information above the next model order. The optimum model order is selected based on testing the highest model orders. A robust, unbiased estimate of minimal rank for information retention providing computational efficiency and improved performance noise discrimination is therewith accomplished.\nU.S. Pat. No. 5,757,675, issued May 26, 1998, to O'Brien, Jr., discloses an improved method for laying out a workspace using the prior art crowding index, PDI, where the average interpoint distance between the personnel and/or equipment to be laid out can be determined. The improvement lies in using the convex hull area of the distribution of points being laid out within the workplace space to calculate the actual crowding index for the workspace. The convex hull area is that area having a boundary line connecting pairs of points being laid out such that no line connecting any pair of points crosses the boundary line. The calculation of the convex hull area is illustrated using Pick's theorem with additional methods using the Surveyor's Area formula and Hero's formula.\nU.S. Pat. No. 6,466,516, issued Oct. 5, 1999, to O'Brien, Jr. et al., discloses a method and apparatus for automatically characterizing the spatial arrangement among the data points of a three-dimensional time series distribution in a data processing system wherein the classification of the time series distribution is required. The method and apparatus utilize grids in Cartesian coordinates to determine (1) the number of cubes in the grids containing at least one input data point of the time series distribution; (2) the expected number of cubes which would contain at least one data point in a random distribution in said grids; and (3) an upper and lower probability of false alarm above and below said expected value utilizing a discrete binomial probability relationship in order to analyze the randomness characteristic of the input time series distribution. A labeling device also is provided to label the time series distribution as either random or nonrandom, and/or random or nonrandom within what probability, prior to its output from the invention to the remainder of the data processing system for further analysis.\nThe above cited art, while extremely useful under certain circumstances, does not provide sufficient flexibility in processing different dimensionalities of data sets of sonar data. Consequently, those of skill in the art will appreciate the present invention which addresses these and other problems."}
{"text": "The present invention is concerned with a decorative method and a blank for a decorator\"\"s tool and in particular a decorative method for providing a grained appearance on a substrate.\nGrain imitation techniques are known in which a scumble glaze is applied to a painted wood substrate, the scumble glaze being applied in a streaky discontinuous manner to give the appearance of natural graining. Currently used techniques however are very time consuming and require considerable expertise, such as that of a craftsman, to create a desired natural grain appearance. They do not therefore appeal to amateur (that is, xe2x80x9cdo-it-yourselfxe2x80x9d) painters, who represent an increasingly important part of the market.\nSimilarly, it is known to create imitation woodgrains on a surface using specialised tools to create the grain finish. Normally, at least two tools are required in order to achieve the desired effect. Existing tools consist of combs having teeth or the like which are used to establish a continuous streaked or grained appearance on a painted substrate, and a separate tool having a convex surface, with a series of concentric ridge formations on the convex surface, which can be used to selectively expose parts of the surface, so as to produce a simulated woodgrain appearance on the surface. Using separate tools to create the woodgrain finish suffers from a number of drawbacks. In addition, packaging of two separate tools in a box or the like requires extra packaging space and packaging material, which adds to the total cost of the product.\nIt is the primary object of the present invention to provide a method which will appeal to the amateur and alleviate some of the abovementioned problems.\nIt is a further object of the invention to provide a blank for a unitary, hand-holdable tool, suitable for use by both amateur and professional decorators in creating an imitation woodgrained finish on a surface.\nAccording to a first aspect of the present invention, there is provided a method of providing a decorative coating for an opaque surface, which method comprises applying to said opaque surface a substantially continuous layer comprising an intimate mixture comprising a particulate (preferably non film-forming) material and an aqueous film-forming polymer binder, said particulate material and said layer both being lighter in color than said opaque surface, and discontinuously removing part of said layer (preferably in streaks) using a tool so as to selectively expose part of said opaque surface.\nAny experienced professional or amateur painter will know, by judgement by eye and instinct, the difference between lighter and darker shades of color. However, if needed, guidance can be gained from the international system of color definition known as the Natural Color System (xe2x80x9cNCS.xe2x80x9d) The NCS system is described in the xe2x80x9cICI Colour Dimensions Colour Atlasxe2x80x9d published by Imperial Chemical Industries P1c of London in 1986 (xe2x80x9cColour Dimensionsxe2x80x9d is a trademark of Imperial Chemical Industries). The NCS system defines color in terms of a cypher, the first two digits of which extend from 00 to 99 with 00 representing white (that is, the ultimate lightness) and 99 representing black, and the intermediate values from 01 to 98 representing increasingly darker shades. Lighter shades therefore have a lower pair of first digits. It is preferred that the particular material and the continuous layer have an NCS value at least 10 units lower than that of the surface.\nThe particulate material is typically an inorganic material, which may be hydrated, such as plaster (gypsum), clay or the like, or non-hydrated, such as chalk or titanium dioxide. When an inorganic material is used, it is advantageous that the intimate mixture used in the method according to the invention may be fire retardant.\nAlternatively, an organic particulate material, such as hollow or alveolate beads of polystyrene or the like, or finely chopped fibers may be used. Suitable organic beads are described, for example, in European Patent Specification No. 0,113,435 B1, published Jun. 7, 1989. This European Patent Specification describes an aqueous paint having an improved xe2x80x9chiding quality.xe2x80x9d The aqueous paint in the European reference consists essentially of pigmented vesiculated beads and opaque polymer particles in a ratio of 30/70 to 90/10, the proportion of opaque polymer particles being smaller at higher pigment volume concentrations. The particulate material is intimately mixed with the aqueous polymer binder, preferably together with thixotropising material, such as a polyurea adduct or bentonite clay, so as to provide thixotropic mixture.\nIt may be preferred that the intimate mixture further comprises reinforcing fibers which can serve to strengthen a resulting coating provided on the opaque surface.\nAdvantageously, the intimate mixture further comprises a drying retardant (a material which retards drying of the coating). The use of a drying retardant is beneficial in allowing the mixture to be applied as a substantially continuous layer and further allowing subsequent removal thereof before substantial drying of the layer has occurred.\nA preferred drying retardant comprises a gel material. Examples of gels suitable for use as a drying retardant are cellulose-based products such as carboxymethyl cellulose, hydroxymethyl cellulose and methyl cellulose, acrylamide and acrylate polymers and copolymers, gelatin, polysaccharides, polyoxamers (polyoxyethylene-polyoxypropylene block copolymers), pectins and agar.\nIt is preferred that the aqueous polymer binder is translucent. Preferably the binder comprises an aqueous dispersion of particles of polymer which are capable of coalescing as the dispersion dries so as to form a film of polymeric material. Suitable polymers include polymers and copolymers of esters such as methyl, ethyl, propyl and hexyl esters of acrylic or methacrylic acids, optionally with acrylic or methacrylic acid or polymers, or copolymers of vinyl esters including vinyl acetate. Certain copolymers of these types are now available as aqueous solutions as opposed to dispersions.\nThe aqueous polymer binder may, in some cases, contain a minor proportion by weight of an organic solvent, which should be miscible with water. An example of such a solvent is an alcohol.\nIt is preferred that the intimate mixture comprises substantially equal amounts of an aqueous carrier, typically water, and the particulate material, the aqueous carrier and particulate material being present in excess of the polymer binder. A preferred intimate mixture comprises 4 to 6 parts by weight of the particulate material, 4 to 6 parts by weight of water optionally 0 to 2 parts organic (preferably alcoholic) co-solvent to assist film formation and 0.5 to 3 parts by weight of the aqueous polymer binder. In a particularly preferred embodiment of the present invention, wherein the intimate mixture further comprises a drying retardant substantially as hereinbefore described, the preferred intimate mixture includes 0.5 to 3 parts by weight of the drying retardant.\nThe opaque surface may be provided by wood, plastics, metal or the like, and it preferably has a darker color than the intimate mixture. In some embodiments it is preferred that the substrate comprises wood. The invention is particularly advantageous in enabling the production of an article having a grained appearance, such as that simulating limed oak, limed maple, limed ash or the like. The substrate may, for example, be timber, or a fiber board such as MDF (medium density fiber board).\nAdvantageously, the substrate comprises a building structure such as a wall panel, door, door frame, window frame, wainscotting or the like. Alternately the substrate may comprise a furniture component such as a fitted kitchen or bedroom unit, closet door, table, seat or the like. In some embodiments, the substrate comprises at least one generally planar surface which is required to have a grained appearance as described above.\nIt is preferred that the method involves applying an opaque protective or decorative coating to the substrate so as to provide the latter with a sufficiently dark opaque surface substantially as hereinbefore described. The opaque coating is preferably applied so as to substantially cover an exposed surface of the substrate. Preferably the opaque coating is water based, although organic solvent-based coatings may be employed.\nIt is further preferred that the method according to the invention involves applying a translucent or transparent coating (such as a varnish, which is preferably water-based), over the discontinuous decorative coating described above. The translucent or transparent coating (which may be lighter or darker than the previously mentioned decorative coating) may itself contain a stain material, pigment, dye or the like; alternately, a stain may be applied to the translucent or transparent coating.\nThe opaque surface typically comprises a paint selected to provide a surface of the substrate with a desired base color. A particularly preferred color of the opaque surface is brown or beige, which can contrast with the lighter color of the intimate mixture, and give the appearance of a naturally grained effect. Further preferred colors of the opaque surface include gray, or other background colors, such as blue, green, red, yellow or the like.\nThe particulate material preferably has a white, cream, beige or brown color which contrasts with the darker color of the opaque surface. Examples of suitable such colors are Y 00 R (usually written as xe2x80x9cYxe2x80x9d) up to Y 90 R according to the NCS definition referred to above. The particulate material may advantageously be mixed with colorants or the like to provide the intimate mixture with a desired color. The color combination of the opaque surface and the intimate mixture is advantageously selected to simulate a grained appearance of a desired wood effect.\nThe decorative coating may be applied by brushing, pad spreading, rolling or spraying onto the opaque surface. Advantageously the decorative coating is allowed to substantially dry (typically for at least one hour) before subsequent application of the intimate mixture thereon.\nThe intimate mixture is preferably similarly applied to the opaque surface by brushing, pad spreading, rolling or spraying thereon. The intimate mixture is preferably applied to provide a substantially continuous layer having a thickness (prior to drying) in the range of 0.3 to 2 mm. The thickness of the layer is selected to provide an appearance of a desired grain, a preferred thickness being about 0.5 mm.\nIt is preferred that removal of the continuous layer involves at least two stages. In the first stage, it is preferred that the continuous layer is discontinuously removed in non-linear streaks to give the opaque surface a veined, streaked or grained appearance, preferably using a comb member including a plurality of distal teeth, the comb member being drawn along the opaque surface to provide a veined appearance resembling that of woodgrain. Preferably the tips of the teeth each respectively comprise a contact edge extending for about 2 to 5 mm in the general direction of the comb member, the tips advantageously being spaced apart by about 2 to 10 mm.\nAdvantageously, in a second stage of the removal of the continuous layer, it is preferred that the veined, grained or streaked surface is treated with at least one tool member having a plurality of spaced apart proud formations (typically of plastics or elastomeric material or the like), to selectively expose parts of the opaque surface.\nIn a highly preferred embodiment, the tool member is a rockable tool, which preferably comprises a convex face, such that the tool can be rocked about the convex face, with a plurality of substantially concentric raised arc formations on the convex face.\nA xe2x80x9crockingxe2x80x9d action is, of course, quite different from a sliding or rolling action. When the tool is rocked, it is caused (generally by a cocking action of the user\"\"s wrist) to pivot, or partially rotate, through no more than 180 degrees, about an axis substantially parallel to the surface, as it is drawn across the surface, before being partially rotated in the opposite sense. This rocking is generally a reciprocal or oscillatory action, in which the tool is rocked to-and-fro in a manner similar to at least one oscillation of a see-saw. This rocking action is generally achieved manually by the user cocking his wrists as the tool is drawn across the surface.\nAs the tool is rocked, the proud formations of the tool are in contact with (and are drawn along) the layer on the substrate. The rockable tool is typically drawn along the opaque surface, advantageously in the general direction of the previously formed veined, grained or streaked appearance, so as to provide a knot or heart grain effect (sometimes known as crown grain effect) on the opaque surface. This knot or heart grain has been found to be especially aesthetically attractive, being a simulation of a cut lengthwise along the trunk of sawn timber (parallel or tangential to growth rings and simulating a cut through a knot region of cut timber).\nIt is further preferred that as the rockable tool is drawn along the opaque surface, the convex face thereof is rocked to and fro substantially in the direction of travel of the tool (that is, about an axis substantially transverse to the direction of travel).\nAccording to a second aspect of the invention, there is, therefore provided a blank for forming a decorator\"\"s tool, which blank can be wrapped around and secured to itself by securing means to form said tool, such that said tool comprises at one end thereof a hollow body having a convex surface with a plurality of spaced-apart proud formations on the convex surface.\nThe proud formations on the convex surface may comprise a series of substantially concentric formations. Advantageously, the concentric formations comprise a series of alternating ridges and furrows concentric about a substantially intermediate point on the convex surface of the tool. The tool can be rocked about the convex surface to complete a simulated woodgrained finish created by the concentric ridges and furrows. Alternately, the proud formations may be more irregular, or may comprise a series of spaced ridges in the form of straight or wavy lines. It may be desirable to provide gaps in such lines at spaced intervals (that is, the lines may be discontinuous).\nThe tool is generally hollow, and is generally also resilient; this means that the user can vary the effect on the surface by varying the pressure applied while drawing and rocking the tool. This represents an operational advantage relative to a non-resilient tool. Preferably, a plurality of apertures are positioned within the alternating ridges and furrows on the convex surface of the tool. Thus, a build-up of excess woodgraining material within the furrows may be substantially alleviated as the material can pass into the interior of the tubular body of the tool, thus facilitating cleaning of the tool and collection of the excess material when such material has passed through the apertures on the convex surface, to the inside of the hollow body of the tool, during the woodgraining process.\nAccording to a preferred feature of the present invention, the blank may contain a plurality of teeth disposed at one end of the tool. Thus, advantageously, in this embodiment both working surfaces needed to create the desired effect are provided in a single hand-held portable unit, which can be formed from an easily storable blank, and either end may be held, as appropriate, when the tool is used to create the grained appearance.\nIn one embodiment of the invention, the teeth may be on a separate attachment which can be secured to the tool. The teeth on the separate attachment may also function as the securing means of the tool which advantageously fits over the ends of the blank when wrapped around itself to form the tool. In another embodiment, the teeth may be on one or both ends of the blank which forms the tool. Where the teeth are on both ends of the blank, two sets of teeth may be provided which advantageously may be orientated in substantially opposite directions on the tool.\nWhen the teeth have been used to produce a veined appearance on the surface, the ridges on the convex surface, when rocked about the convex surface, complete the continuous natural grained finish. Typically, the tips of the teeth comprise a contact edge in which one face is relatively more bevelled (a so-called xe2x80x9cchiselxe2x80x9d type tooth), or the two faces may be at substantially the same angle of inclination (a so-called spear-section). The respective inclined faces may typically extend for about 2 to 5 mm from each tip of the teeth.\nIn some embodiments, either side of the teeth can be used in either direction, eliminating confusion as to which side of the tool is to be used. Typically, the teeth and ridge formations are from 2 to 10 mm apart.\nIn a preferred embodiment, the teeth, the convex surface and the proud formations all together comprise a unitary integral molding.\nIn another embodiment, the proud formations may be on a separate molding attachment, releasably engageable with the blank which comprises the base of the tool. This would allow imitation woodgrains of different size and shape to be created by the same tool.\nTypically, the securing means comprise complementary jointing formations such as a dovetail joint, or pop studs releasably engageable with reception apertures on the blank. In the latter case, further reception apertures may be provided in series along the length of the blank to receive such pop studs, serving to maintain the blank in its proper configuration. Alternate securing means comprise adhesive (for example double-sided adhesive tape). Thus, advantageously, the diameter of the convex surface of the tool may be varied, thus creating grains of different size and shape.\nConnection means may also be provided to connect the molding attachment to the tool; the connection means may comprise pop studs releasably engageable with reception apertures of the type as described for maintaining the blank in the tool configuration.\nThe blank may be made of any suitably stiff material, but which is sufficiently flexible to allow the blank to be wrapped around, or doubled over on itself, without snapping. The material used for the blank may comprise a suitable plastics material, which is sufficiently flexible and lightweight, and thus easily portable. As mentioned above, the resulting hollow tool is preferably resilient.\nIn one embodiment of the invention, the blank is substantially flat, and may have a plurality of fold lines along its length such that when the blank is wrapped around itself along the fold lines it forms the tool. Alternatively, the blank may be stepped in cross section.\nThe present invention is particularly advantageous in providing a decorated article, and a method of producing the same, which are largely free of organic solvents (except possibly for small amounts of co-solvent). This aspect is environmentally beneficial, which is particularly advantageous when the resulting articles are found or used within building interiors.\nThere is further provided by the present invention an article having a decorative coating, and obtained according to a method substantially as hereinbefore described.\nThe method according to the present invention results in an article comprising:\n(a) a substrate having an opaque surface;\n(b) a decorative coating provided on said surface and comprising an intimate mixture comprising a particulate material and an aqueous film-forming polymer binder, said coating being lighter in color than said surface and being discontinuously provided thereon so as to expose part of said surface.\nThe nature of the substrate, opaque surface and decorative coating are substantially as hereinbefore described."}
{"text": "This invention relates to computer programming, and more particularly to computer aided software development.\nThe purpose of the invention is to provide a programming environment designed to enhance the speed and productivity of software development, particularly a method for substantially decreasing the time required for recompilation and relinking in the edit-compile-link-run cycle of the software development process. When code is being written, the elapsed time through the edit-compile-link-run cycle after the user makes a small change to the application source code is called the turnaround tune. A primary purpose of the invention is to minimize this turnaround time.\nThe programming \"environment\" as the term is used herein means the set of programs or modules (i.e., code) used to implement the edit-compile-link-run cycle for a developer, who is ordinarily seated at a terminal and engaged in the endeavor of writing code. The environment which is the subject of this invention will be called \"the environment\", whereas any program being developed under the environment will be called \"an application\". The environment is capable of supporting the development of any application, including the environment itself. The user of the environment is called \"the developer\", while the user of an application is called \"the end user\".\nSoftware development is characterized by a process involving the steps of editing the program, compiling and linking the program, and running the program. A compiler translates a source program that has been written in a high-level language such as Pascal or Fortran into a machine executable form known as an object program.\nThe software development process is further divided into stages, with the earlier stages characterized by rapid and large scale activity (e.g., editing) in all or most of the application source files, and in the later stages characterized by less frequent and smaller changes in fewer than all of the source files. During the earlier stages the objective is removing syntax errors in the source code and logic errors in the application. During the final stages the objective is improving the efficiency of the application and testing the behavior of the application in the form it will be delivered to the end user.\nIt is generally desirable that the quality of the object code generated by a compiler, as measured in terms of efficiency, be as good as possible. A compiler that generates very efficient object code is known as an optimizing compiler. Optimized object code is characterized by maximized efficiency and minimized execution time. However, the complex methods and techniques employed by optimizing compilers to produce highly efficient object code necessarily result in relatively long compile times.\nThe removal of logic errors is relatively independent of the efficiency of the implementation of the application; therefore, during the early stages of software development, it is desirable that the environment emphasize turnaround time over optimization. In addition, during the early stages it is advantageous to insert application-run-time checks for certain kinds of detectable faults such as boundary overrun. The concerns with efficiency and testing during the final stages require optimization, and the lower frequency of changes makes the use of traditional software tools effective.\nThe edit-compile-link-run cycle is typically repeated numerous times during development of a particular piece of software. At any stage of this activity the developer may be required to correct detected errors (as used herein, \"error\" means \"need for a change\", since the motivation for making a change may be either repairing a previous oversight or adding new functionality). Errors may be detected by the compiler, the linker, or later by the programmer during test execution. This style of interaction in conventional environments results in frequent context changes and delays for the developer. Context changes occur while the developer separately and sequentially invokes the editor, the compiler, the linker, and the application itself. Delays occur while the developer waits for these separate tools to complete their tasks.\nThus, while long compile times are tolerable in the final stages of developing an application, i.e., when generating production quality object code, these delays are not tolerable in the early stages of the process of developing, testing, and debugging software where long compile and link times will be much more noticeable since these are invoked much more often. Moreover, changes in the application code made during development are usually localized and small in size with respect to the rest of the program. In known software development environments, the turnaround time for compiling an application module is proportional to the size of the module and the turnaround tune for linking an application is proportional to the size and number of modules. In the environment of this invention both compilation and linking turnaround are proportional to the size of the changes to the source code made by the developer since the last compile/link operation. Many applications programs have 100,000 to 1,000,000 or more lines of code; the turnaround time (time for the edit-compile-link-run loop to complete) in developing such programs can become an overnight activity, and thus presents a major burden.\nThus, it is desirable to provide a software development environment that would allow fast turnaround in the edit-compile-link-run cycle.\nExamples of commercially available developmental compilers include \"Quick C.TM.\" by Microsoft Corporation, \"LightSpeed C\" by Symantek Corporation, \"Turbo C\" by Borland Corporation, and Saber-C by the Saber Company. These prior systems are faster than traditional batch compilers, and may provide some degree of incremental (as distinguished from batch) operation; for example, a module may be treated separately if only that module has been changed since the last edit-compile-link-run cycle. This level of incremental operation is known as coarse-grained incremental operation."}
{"text": "1. Field of the Invention\nThis invention relates to a multiple input restrictive logic OR circuit and more particularly to a high speed fully differential logic OR circuit with multiple non-overlapping inputs.\n2. Description of the Related Art\nWhen designing systems such as automated test equipment it is desirable to provide a high speed logic circuit which accepts numerous inputs and provides a single output, where only one input at a time is active and the output is active when one of the inputs is active. In such high speed applications, input-to-output propagation delays must be minimized while remaining uniform for each input. It is also desirable to minimize signal noise and interaction between circuit devices that may degrade the signal.\nDifferential logic OR circuits have been developed using known ECL technology such as the differential two input Motorola MC10EL05. However, these circuits commonly have only 2 or 3 inputs which suffer from different input-to-output propagation delays. In addition, to construct a logic OR circuit with more inputs, for example 8, several 2 or 3 input OR circuits must be coupled together in a tree structure. Such tree structures have different input to output propagation delays for the different inputs depending on the path taken through the logic tree. They also result in undesirable interaction between the multiple logic gates, which can degrade the signals.\nA multi-input logic OR circuit could also be implemented with a \"Wired-OR\" circuit, in which single-ended input lines are combined to produce an \"OR\" output. An example of this approach is embodied in the Motorola MECL 10K 4 input wired-OR logic gate. This circuit has relatively low propagation delay and is generally symmetric with respect to propagation delay between input and output. However, its single-ended inputs make the circuit susceptible to jitter and makes propagation delay sensitive to environmental changes.\nA multi-input logic OR circuit could also be implemented with a multiplexer that selects data from one of the multiple input lines and directs the data to a single output line. An example of this approach is embodied in the Motorola MC10H164 8-line multiplexer. The selected input line is activated by control lines and the number of control lines increases as the number of inputs increase. A four input multiplexer requires two control lines, while an eight input multiplexer requires three control lines. These circuits suffer from undesirable input-to-output propagation delays and experience an additional delay in the settling time of the control lines. The necessity for control lines makes this approach significantly more complex."}
{"text": "Storing and safeguarding electronic content is of paramount importance in modern business. Accordingly, large storage systems may be utilized to protect such electronic content. In order to enhance the delivery of such electronic content to requesting devices, such large storage systems may utilize journaling systems. Unfortunately, the use of multi-core microprocessors may complicate the use of such journaling systems due to their ability to parallel process multiple write requests."}
{"text": "1. Field of the Art\nThe present invention relates in general to a hydraulically-operated anti-skid brake system for a motor vehicle, and more particularly to improvements in consistency of braking effect of right and left wheel brake cylinders upon anti-skid brake application while right and left wheels are running on road surfaces of different conditions.\n2. Related Art Statement\nSuch an anti-skid brake system uses, for example, an anti-skid device of a so-called closed-circuit type or variable-volume type as disclosed in Japanese Patent Application which was laid open under Publication No. 58-26659, or an anti-skid device of a so-called recirculating type as disclosed in Japanese Patent Application which was laid open under Publication No. 57-80956. In the closed-circuit type, a variable-volume chamber is provided in communication with a part of the fluid passage which communicates with a wheel brake cylinder, but is disconnected from a master cylinder of the brake system upon activation of the anti-skid device. The pressure in the above-indicated part of the fluid passage is regulated by changing the volume of the variable-volume chamber by an actuator. This actuator is operated by a pressure generated by a hydraulic pump, a vacuum pump or other pressure source provided independently of the master cylinder. In the recirculating type, the brake fluid is discharged from a wheel brake cylinder to lower the pressure in the brake cylinder upon an excessive rise in the braking pressure therein, and the discharged fluid is pressurized to a higher level of pressure by a pump. The thus reserved pressurized brake fluid is recirculated into the brake wheel cylinder when it becomes necessary to raise the braking pressure in that brake cylinder. In either one of the above two different types, solenoid-operated control valves are generally provided for the right and left wheel cylinders, and are controlled independently of each other by a main controller whose major part consists of a computer. For example, the main controller is adapted to receive signals from slip sensors which detect slippage of the right and left wheels of the vehicle on the road surfaces. Based on these signals, the main controller controls the solenoid-operated control valves so as to hold the slip ratios of the right and left wheels within an optimum range.\n3. Problem solved by the invention\nHowever, the hydraulic brake system equipped with an anti-skid device indicated above suffers from a problem of reduced straight-line running stability of the vehicle in the event that relatively high braking pressures are applied to the brake cylinders while the right and left wheels are running on road surfaces having different coefficients of friction. Described more specifically, the solenoid-operated control valve corresponding to the wheel running on the road surface with a comparatively low coefficient of friction starts an anti-skid control of the braking pressure in the corresponding brake cylinder, at an earlier timing than the other control valve, so as to hold the braking pressure in that brake cylinder at a comparatively low level. On the other hand, the other control valve corresponding to the other wheel running on the road surface with a higher coefficient of friction is left in its non-operated position, allowing the braking pressure in the corresponding brake cylinder to be raised at a higher rate. As a result, the right and left brake cylinders have different braking effects on the corresponding wheels of the vehicle, whereby the straight line stability of the vehicle is reduced."}
{"text": "Video conferencing is becoming an important way for people and businesses to communicate. However, challenges to mobile video conferencing remain, for example, many devices require the user to hold their phone in front of their face to capture an image of their face. What is needed are devices and systems that allow hands-free mobile video conferencing."}
{"text": "Cameras are used to capture photos. Traditional cameras capture photos on film. The film is developed and the photos are printed. The negative, i.e. the original developed film, provides a source from which new prints can be obtained. The advent of high-quality scanning also allows a print to act as a source of copies, although not with the same quality as from a negative. If either a negative or a print is damaged, the quality of a subsequent copy is compromised.\nDigital cameras capture photos as digital images which can be both displayed and printed. A digital image thus acts as a digital “negative”. However, many users find the task of printing and otherwise managing digital images onerous."}
{"text": "Conventionally, the material forming an electric wire conductor used to be routed in an automobile or the like is mostly a copper-based material such as copper or copper alloy which is excellent in electrical conductivity. Various studies have been conducted for improving mechanical properties of the conductor such as tensile strength (see for example PTD 1 and PTD 2).\nPTD 1 discloses an electric wire conductor for an automobile. This conductor is made up of a plurality of stranded hard wires made of a copper alloy. The copper alloy contains any one element selected from Mg, Ag, Sn, and Zn at a content in a specific range. The copper alloy is subjected to wiredrawing at a degree of cold working of 99% or more to thereby enhance mechanical properties such as tensile strength, Young's modulus, and electrical conductivity. PTD 2 discloses an electric wire for a wire harness. The electric wire includes a conductor of a copper alloy. The copper alloy contains Ti at a content in a specific range as a precipitation strengthening element and Fe at a content in a specific range as a precipitation promoting element. Thus, additive elements dissolved in a Cu matrix are effectively precipitated, and accordingly mechanical properties such as tensile strength and electrical conductivity are enhanced."}
{"text": "Distributed denial of service (DDoS) attacks present security and availability issues for many organizations, and in particular, for enterprises engaged in content delivery services. In a DDOS attack, many distributed hosts flood a target system with traffic, such as HTTP requests directed at a web server under attack. The flood of traffic overloads the server so that the system under attack cannot respond to legitimate traffic in an effective manner. Such attacks, and the resultant unavailability, can produce several adverse consequences for the operator of the server, including loss of reputation, potential loss of business or revenue, and substantial bandwidth costs.\nA conventional technique for responding to DDoS attacks is to use customer-triggered real-time black holes, otherwise known as remote triggered destination internet protocol (IP) address black hole filtering (RTDBHF). Black hole filtering (BHF) results in packets being forwarded to a router's bit bucket (e.g. Null 0/discard interface/null interface). Traditionally, RTDBHF works solely based on the destination address of the traffic by exploiting the forwarding logic of routers. All traffic to the attacked DNS or IP address is sent to the null interface. RTDBHF allows destination IP address black holes to be triggered remotely, by customers, or an internet service provider. A user can remotely trigger a destination address network-wide black hole filtering response using border gateway protocol (BGP) and static routes pointing to the null interface. Thus, although RTDBHF discards attack traffic directed towards the destination and mitigates collateral damage to other systems and network availability, the targeted system is taken completely offline as both legitimate traffic and attack traffic to the destination address are discarded.\nDestination IP address enhanced BGP-triggered black holing techniques, also known as remote triggered destination enhanced black hole filtering (RTDEBHF) have been developed that address this concern. RTDEBHF techniques uniquely identify autonomous system (AS) border routers that could direct attack traffic to the targeted system. BGP community values are also assigned to identify sets of the border routers. By using a customized internal BGP (iBGP) advertisement containing the address of the targeted network and BHP community value, only the next hops of the selected routers are changed to the null interface, and the original next hop addresses to the targeted network on all other routers are preserved. Thus, traffic is filtered only from the routers identified as routers that could direct attack traffic and having specific route map matches for the BGP community value, while all other traffic will get forwarded to the targeted network.\nAn alternative conventional technique for handling DDoS attacks is to use remote triggered black hole filtering with Unicast Reverse Path Forwarding (uRPF), also known as remote triggered IP address source black hole filtering (RTSBHF). RTSBHF is a technique that allows black hole filtering based on the source address of the network traffic. uRPF techniques are combined with remotely triggered black hole filtering so that BGP can be used to distribute discard routes directed to the null interface, based on the source address of the attack traffic. This results in all traffic to and from a source address to be dropped.\nUp until now combining RTDBHF with RTSBHF has not been considered feasible. Thus, there is a need for more robust solutions to provide the combined benefits of RTDBHF and RTSBHF."}
{"text": "Downhole constructions including oil and natural gas wells, CO2 sequestration boreholes, etc. often utilize borehole components or tools that, due to their function, are only required to have limited service lives that are considerably less than the service life of the well. After a component or tool service function is complete, it must be removed or disposed of in order to recover the original size of the fluid pathway for uses such as hydrocarbon production and CO2 sequestration. Disposal of components or tools can be accomplished by milling or drilling the component or by tripping the tool out of the borehole. Each of these is generally time consuming and expensive. The industry would be receptive to new materials, and methods that remove a component or tool from a borehole without such milling and drilling operations."}
{"text": "(a) Technical Field\nThe present invention relates to a system for charging an electric vehicle using an in-cable control box (ICCB).\n(b) Description of the Related Art\nRecently, a system for slowly charging an electric vehicle has been used, in particular, electric vehicle station equipment (EVSE), and an electric charging system using an in-cable control box (ICCB).\nAn electric charging system using the EVSE is used by installing the EVSE, which is connected to an electricity distribution board, at a necessary position for performing work for supplying electric power. According to the electric charging system using the EVSE, the time required to charge the electric vehicle is comparatively short, and the time required to charge the electric vehicle may be reduced because of a capacity of the EVSE. In particular, an electricity rate applied for charging the electric vehicle, which is set by an electric power company, is the same for all users, such that a progressive tax is not imposed, and thus the electricity rate is low.\nReferring to FIG. 2 (PRIOR ART), the electric charging system using the ICCB is used to charge the electric vehicle by connecting the ICCB 2 to a general power socket 1, which is used in the related art, and has an advantage in that the electric vehicle may be charged at any place where the power socket is installed, and separate installation costs are not required.\nHowever, in the case of the electric charging system using the EVSE, there are problems in that high installation costs are incurred due to costs required to purchase the EVSE and costs required for work for supplying electric power, and there is a limitation on the type of place for charging the electric vehicle.\nThere are problems in that in the case of the electric charging system using the ICCB, because the electric vehicle is charged using a general domestic power socket, it is impossible to distinguish whether to charge the electric vehicle or a general electric device, and therefore, as electric consumption increases, a progressive tax is imposed, such that excessive costs are required for maintenance of the vehicle, and time required to charge the electric vehicle is increased in accordance with the rating of the power socket connected with the ICCB.\nThe above information disclosed in this Background section is only for enhancement of understanding the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "1. Technical Field\nThe present invention relates in general to read and write commands to nonvolatile memory devices within a data processing system and in particular to read and write commands from a controller to a Redundant Array of Independent Disks within a data processing system. Still more particularly, the present invention relates to reducing the number of parity read and write commands between a controller and a Redundant Array of Independent Disks, Level 6.\n2. Description of the Related Art\nA Redundant Array of Independent Disks (\"RAID\") is an array, or group, of hard disk drives controlled by a single array controller and combined to achieve higher transfer rates than a single, large drive. Even though multiple drives are controlled by one adapter, the RAID device appears as one drive to the data processing system. Depending on the configuration, the RAID device will increase the level of protection and storage capacity for a data processing system over a single, hard disk drive. The primary functions of the RAID system are to increase the availability, protection and storage capacity of data for a data processing system.\nRAID technology generally distributes data across the drives according to the format of the particular RAID classification (RAID 1, 2, 3, 4, 5 or 6). Copies or portions of data for a particular file may be written in segments on more than one disk drive, a process referred to as \"striping.\" By storing the data and instructions on multiple drives, higher data transfer rates are enhanced by the ability of the controller to schedule read and write commands to multiple drives in parallel.\nRAID 5 reads and writes data segments across multiple data drives and writes parity to the same data disks. The parity data is never stored on the same drive as the data it protects, allowing for concurrent read and write operations. Within any stripe of a five drive RAID 5 configuration, all drives contain data information and parity information. If one of the data drives were to fail, the remaining four data drives and the parity on each remaining may be used to regenerate user data which improves improving data protection.\nRAID 6 improves the data protection of RAID 5 by providing two parity drives. The original technique for data protection in RAID 6 was to copy the parity drive onto a parallel parity drive, or \"mirror\" the parity drive. This protects the RAID 6 device from a parity drive failure, but does not protect the group from failure of two data drives. In order to protect against multiple data drive failures, RAID 6 changes the configuration so that the second parity drive will protect across different drive groups. For instance, parity drives are arranged so that each data drive has parity stored on two parity drives. A RAID 6 device with this configuration would be depicted as having multiple rows and multiple columns of data drives with each row and column ending with a parity drive. Parity of each data drive would then be stored on two drives.\nIn large arrays the increase in the number of additional drives is substantial, but not prohibitive. If the array is a ten by ten array, of 120 drives only 20 are parity drives. However, in small arrays the percentage increase is large. For example, if the RAID subsystem contained a four by four array of data drives, a parity drive would be added for each row of data drives. In addition, parity drives would be added for each column. Therefore, for an array containing sixteen data drives, there would be eight parity drives--a fifty percent increase in the number of drives.\nA RAID 6 device provides extra data protection but, at a somewhat prohibitive cost. The two group version of RAID 6 requires that a single data disk belong to two parity groups. If a data drive stripe were to be updated, the parity information of all parity drives affected would also need to be updated resulting in many more reads and writes.\nIn FIG. 5, a RAID 6 configuration is depicted, showing the first several data stripes on the individual drives. The letter \"D\" in the diagram indicates that DATA is stored in that location and \"P\" indicates that PARITY is stored in that location. The number indicates the data segment stored in that location. For example, to calculate the parity information stored in P05, a RAID 6 device would need to read D01, D02, D03, and D04. It would then calculate the parity and write the results to P05.\nIn the two group version of RAID 6, a single disk belongs to two parity groups. In this instance, D01 belongs to a horizontal parity stripe and a vertical parity stripe. In order to update D01, parity information stored at P05 and P21 also needs to be updated.\nUsing the two parity group version, a RAID 6 device could handle up to three drive failures without losing any information. For example, if D01 failed, its information could be retrieved using either the horizontal rank or vertical rank parity drives. If P21 also failed, the vertical rank would not contain enough information to regenerate D01. However, the horizontal rank would be available to regenerate D01. On the other hand, if on the horizontal rank, P05, were to fail then the vertical rank parity drive could be used to regenerate D01. If both the vertical rank, P21, and the horizontal rank, P05, were to fail before D01 were regenerated, then D01 could not be regenerated from either the horizontal or vertical ranks.\nA significant problem with RAID 6 devices is the number of parity updates that must be generated. When a data drive is updated, parity needs to be calculated for two drives. This procedure, referring again to FIG. 5, requires a read to the data drive D01 and the parity drives P05 and P21. In addition, a write is required to the data drive D01 and both parity drives, P05 and P21. This is for a single drive stripe write.\nAn example of a full stripe write, in a RAID 6 configuration, is depicted in FIG. 6 (assuming the four by four disk array in FIG. 5). The process begins with step 600, which depicts the host sending a write command to the RAID 6 controller. The process passes to step 602, which illustrates the controller sending a read command to D01, D02, D03, D04 and Parity drives P21, P22, P23, and P24. The process proceeds to step 604, which depicts the controller XORing the old data, the new data and the old parity. This new parity is then written to the parity drives, P21, P22, P23, and P24, in segments across the vertical parity drive stripe. The process then continues to step 606, which illustrates the controller sending a write command to the data drives DO, D02, D03, and D04 and the parity drives P05, P21, P22, P23, and P24. The process then passes to step 608, which depicts the controller sending a completion signal to the host.\nIn summary, the write command requires a read to segments on each of four data drives, e.g., R(1,2,3,4), where R(1,2,3,4) are reads to segments DO, D02, D03 and D04. In addition, reads to the Parity drives R(21,22,23,24) are required for vertical parity calculations, writes W(1,2,3,4,5,21,22,23,24) are needed to write out the user data and all of the parity data (referring to FIG. 5). A read to each of four parity drives and a write to each of the parity drives and data drives are required. A minimum full stripe write to a four by four disk array requires eight Read operations and 9 Write operations, even where the RAID 6 subsystem leverages its horizontal parity by eliminating the need to write to the horizontal parity drive.\nIt would be desirable, therefore, to provide a method for reducing the number of parity writes for RAID 6 devices.\nIt would also be desirable, to reduce the Input/Output load on the controller which will provide room for additional drives in the RAID 6 device controller."}
{"text": "In the field of graphic arts, a printing plate is produced using a set of color separations of a color original prepared using lithographic films. In general, colorproofs are prepared from color separations in order to inspect for errors in color separation and to check the need for color correction and the like before printing. Color proofs are required to realize high resolution enabling accurate half tone reproduction and high processing stability. To obtain color proofs close to actual prints, it is desirable for the materials of color proofs to be the same as those used on press, i.e., the same paper and the same pigments. There is a higher demand for a dry process involving no processing solutions for the preparation of color proofs.\nWith the recent spread of computerized systems in prepress work, recording systems for preparing color proofs directly from digital signals have been developed. Such computerized systems, particularly contemplated for preparing high quality color proofs, are generally capable of reproducing dot images at 150 lines or more per inch. In order to obtain high quality proofs from digital signals, a laser beam is used as a recording head, which is capable of modulation according to digital signals and focusing into a small spot diameter. Hence it is demanded to develop image forming elements that exhibit high sensitivity to laser light and high resolution enabling reproduction of highly precise dot images.\nImage forming elements known useful in laser transfer methods include a thermal melt transfer sheet, which comprises a substrate, a light-heat conversion layer capable of absorbing laser light to generate heat, and an image forming layer having a pigment dispersed in a heat fusible matrix (e.g., a wax or a binder) in the order described, as disclosed in JP-A-5-58045. A thermal transfer sheet of this type is brought into contact with an image receiving sheet and imagewise irradiated with a laser beam. The irradiated area of the light-heat conversion layer generates heat to melt the image forming layer, and the molten part of the image forming layer is transferred to the image receiving sheet.\nJP-A-6-219052 teaches a thermal transfer sheet comprising a substrate, a light-heat conversion layer containing a light-heat converting substance, a release layer as thin as 0.03 to 0.3 μm, and an image forming layer containing a colorant. According to this technique, the release layer reduces its bonding strength between the image forming layer and the light-heat conversion layer upon being irradiated with laser light. As a result, the image forming layer is allowed to be transferred to an image receiving sheet that has been brought into contact with the thermal transfer sheet to form a high precision transfer image. This image formation method utilizes laser ablation. That is, a laser-irradiated part of the release layer decomposes and vaporizes, resulting in reduction of the strength bonding the image forming layer and the light-heat conversion layer in that area. As a result, the image forming layer of that area is transferred to the image receiving sheet.\nThese imaging methods are advantageous in that images can be formed on printing paper having an image receiving layer (adhesive layer) and that a multicolor image can easily be obtained by successively transferring images of different colors onto the same image receiving sheet. The method utilizing ablation is particularly advantageous for ease of forming a highly precise image and is useful to prepare color proofs (DDCP) or precise mask images.\nWith the spread of desk-top publishing (DPT) work, printing companies adopting a computer-to-plate (CTP) system have a strong demand for a DDCP system, which eliminates the need of intermediate film or plate output as has been involved in traditional analog proofing. In recent years, DDCPs with higher quality, higher stability, and larger sizes have been demanded as good approximations to the final prints.\nLaser thermal transfer systems are capable of printing at high resolution. Options include laser sublimation, laser ablation, and laser melt, each of which has the problem that the recorded dot shape is not sharp enough. The laser sublimation system is insufficient in approximation in color to the final print result because of use of dyes as coloring matter. Besides, this system involving dye sublimation results in blurred dot outlines, failing to achieve sufficiently high resolution. The laser ablation system, which uses pigments as coloring matter, provides a satisfactory approximation in color to the final printed products, but the dots are blurred, resulting in insufficient resolution similarly to the dye sublimation system because of the involvement of coloring matter scattering. The laser melt system also fails to create clear dot outlines because the molten colorant flows.\nThe colors that can be reproduced by conventional heat transfer sheets are limited to process colors (i.e., yellow, magenta, cyan, and black) and their combinations. There are needs for metallic spot colors that are difficult to reproduce with the four process colors, such as metallic gold and metallic silver. While it is known that these metallic colors can be reproduced by using tabular inorganic particles, there is still room for improvement in sensitivity and resolution."}
{"text": "In the past, various condition responsive mechanisms utilized a sensing device for ascertaining temperature changes and effecting operation of such condition responsive mechanisms in response to such temperature changes. These sensing devices may comprise an expansible bellows communicated with a generally elongate metallic capillary tube, and such bellows and tube are charged with a temperature sensitive pressure fluid, such as a liquid or a gas for instance. One of the aforementioned condition responsive mechanisms is illustrated in U.S. Pat. No. 3,648,214 issued Mar. 7, 1972 to John L. Slonneger which is incorporated by reference herein.\nNumerous arrangements were known in the past for crimping, i.e. closing and sealing, one or both ends of the aforementioned capillary tube to confine the pressure fluid therein. One typical arrangement was to flatten the tube in a limited region thereof and thereafter effect a welding operation on such flattened region. For example, two or more welds were made across the flattened region of the capillary tube, i.e. transverse to the lengthwise direction of the tube, and the tube was then severed between a pair of those cross-wise welds so as to provide adequate sealed ends of the tube. In another past crimping arrangement, a tube end was crimped into a semi-circular congifuration and thereafter such crimped tube end was sealed by welding or soldering. It is believed that, at least with capillary tube metals, the heating effected during a welding or soldering operation might have a deleterious affect on the physical properties of such metals. In another past crimping arrangement, a metallic capillary tube had its opposite sides collapsed inwardly between a pair of crimping jaws, and while confining these collapsed opposite sides, another pair of crimping jaws were operated to collapse a portion of the tube with sufficient force to sever the tube forming a tapering cold welded end thereon."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an infrared gas analyzer for detecting gases contained in air, exhaust gas, and the like. More specifically, it relates to an improved infrared gas analyzer that prevents influenced results by interfering gases.\n2. Description of Related Art\nAs to an infrared gas analyzer for detecting gases contained in air, exhaust gas, and the like, in the case where interfering component gases having an absorption wavelength range partially overlapping a wavelength range of a gas to be measured are contained in a sample gas, an infrared gas analyzer capable of preventing an influence by the interfering component gases comprising an interference filter as shown in FIG. 2 has been known.\nThis gas analyzer comprises a reference cell 1, in which a reference gas is enclosed. A sample cell 2, in which a sample gas is to be supplied, is disposed in parallel. The sample cell 2 is provided with a supply port 3a for supplying the sample gas and an exhaust port 3b for exhausting the sample gas. Reference numerals 4a, 4b designates a light source disposed on one side of the reference cell 1 and the sample cell 2, respectively. Reference numeral 5 designates a detector, such as condenser microphone, disposed on the other side of the reference cell 1 and the sample cell 2.\nReference numeral 6 designates a chopper disposed between the reference cell 1 and the sample cell 2 and the light sources 4a, 4b. Reference numerals 7a, 7b designates interference filters disposed between the reference cell 1 and the sample cell 2 and the detector 5. The filters 7a, 7b pass radiation including a wavelength range in any absorption bands of the gas to be measured but reflect radiation that is in a non-absorbing wavelength range. Reference numeral 9 designates a preamplifier.\nWith this infrared gas analyzer, a radiation emitted from the light sources 4a, 4b is intermittently incident upon the reference cell 1 and the sample cell 2 by rotating the chopper 6. Thereupon, a part of the radiation is absorbed by a gas enclosed in the reference cell 1 and the sample cell 2, respectively, so that the energy of radiation incident upon the detector 5 from the reference cell 1 is different from that from the sample cell 2. The sample gas is analyzed on the basis of this difference in energy of radiation.\nAnd, radiation including any wavelengths within the absorption bands of the interfering gases contained in the sample gas supplied in the sample cell 2 and incident upon the reference cell 1 and the sample cell 2 are reflected by the interference filter 7 to prevent them from being incident upon the detector 5, whereby the influences of the interfering component gases are reduced.\nAlso an infrared gas analyzer shown in FIG. 3 has been known. In this infrared gas analyzer, the interference filter 7 in the gas analyzer shown in FIG. 2 is replaced with gas filter cells 8a, 8b for absorbing the interfering component gases. Other constructions of FIG. 3 are the same as in the gas analyzer shown in FIG. 2, so that they are marked with the same reference numerals and marks.\nIn the gas analysis by this gas analyzer (FIG. 3), radiation having wavelengths of the absorption bands of the interfering components contained in the sample gas and incident upon the reference cell 1 and the sample cell 2 are absorbed by the gas filter cell 8 to reduce the influences of the interfering components.\nAlso, an infrared gas analyzer of an interference compensation type as shown in FIG. 4 has been known. In this gas analyzer, a detector 5a is adapted to be able to pass radiation therethrough. An interference-compensating detector 5b, upon which the radiation passing through the detector 5a is incident, is provided. A subtracter 10 for subtracting an output of the interference-compensating detector 5b from an output of the detector 5a is provided in the gas analyzer as shown in FIG. 4.\nOther constructions are the same as in the gas analyzer shown in FIG. 2, so that they are marked with the same reference numerals and marks.\nWith the conventional infrared gas analyzer as shown in FIG. 2, a reflection factor of the interfering component of radiation by the interference filter 7 is high, so that the influences of the interfering component radiation can be reduced.\nHere, the transmission and reflection of the radiation by the inteference filter in the case where the window on the side of the detector of the gas cell S is replaced with the interference filter f, as shown in FIG. 5, is investigated.\nRadiation which shall be illustrated by a light beam I.sub.o emitted from a light source L passes through the window W of the gas cell S to enter the gas cell S. After a part of the light I.sub.o is absorbed by the component to be measured contained in the gas cell S, the remaining light I.sub.o passes through the interference filter F to enter the detector (not shown).\nThe total transmission quantity T.sub.1 of radiation having a wavelength range capable of passing through the interference filter F at this time is expressed by the following equation (1): ##EQU1## Even though tf is replaced with tw and rf is replaced with rw, it is one and the same thing, so that the transmission quantity of radiation capable of passing through the interference filter has nothing to do with their incident direction. That is to say, it is one and the same thing even though the radiation is are incident from the side of the interference filter F.\nOn the other side, the total reflection quantity R of radiation having a reflection wavelength range of the interference filter is expressed by the following equation (2): ##EQU2## The quantity of the reflected radiation is changed by replacing tf with tw and rf with rw. That is to say, the case where the radiation is incident from the side of the interference filter F is different from the case where the radiation is incident from the side of the window of the gas cell S.\nIt is found from the above described matters that since both the window W of the gas cell S and the interference filter F have a high transmittance and a low reflectance for the radiation of a component to be measured, the quantity of the reflected radiation can be deemed as constant regardless of the direction of the surface upon which the radiation is incident. However, the filter has a low transmittance and a high reflectance for the interfering component radiation, so that the quantity of the radiation reflected toward the side of the radiation source in the case where the radiation is incident from the side of the filter is larger than that in the case where the radiation is incident from the side of the window W of the gas cell S.\nIf the radiation reflected toward the side of the radiation source is not return by reflecting again by means of a radiation source mirror and the like, there is no difference in quantity of the transmitted radiation regardless of the direction of the surface upon the light is incident.\nHowever, in order to increase the quantity of radiation in fact, a mirror is frequently used with the light source, so that the difference in incident direction always leads to a difference in quantity of transmitted radiation.\nI.sub.o designates a quantity of an incident radiation; tw designates a quantity of a transmitted radiation; rw designates a quantity of a reflected radiation by the window W; tf designates a quantity of lights passing through the interference filter F; rf designates a quantity of a reflected radiation by the interference filter F; c designates a concentration of a gas contained in the cell S; l designates a length of the cell S; and .epsilon. designates a constant determined by the gas.\nAccordingly, referring to FIG. 2 and FIG. 4, a comparatively large quantity of interfering component radiation reflected by the interference filters 7a, 7b is reflected by the inside surface of the reference cell 1 or the sample cell 2, the windows of the reference cell 1 and the sample cell 2, reflecting mirrors of the light sources 4a, 4b and the like to arrive at the interference filters 7a, 7b again. This is repeated.\nMoreover, since the interference filters 7a, 7b are disposed between the reference cell 1 and the sample cell 2 and the detector 5, a large quantity of the interfering component radiation is reflected by the inside surface of the reference cell 1 and the sample cell 2 to increase an oblique component incident upon the interference filters 7a, 7b. In general, a transmission spectrum of the interference filter has physical characteristics of shifting toward shorter wavelengths in the case of an oblique incidence. Accordingly, as a result, a problem occurs in that the quantity of the interfering component radiation passing through the interference filters 7a, 7b to enter the detector 5 is increased which lowers the accuracy of analysis.\nNext, an infrared gas analyzer shown in FIG. 3 has no problem in the reflection of the interfering component radiation in the interference filter since the radiation having wavelength within the absorption bands of the interfering component gases are absorbed by the gas filter 8.\nHowever, in the case where the degree of absorption of the interfering component radiation is slightly lower and a plurality of kinds of interfering component gases are contained in the sample gas, it is difficult to absorb all interfering component radiation, so that the quantity of the interfering component radiation incident upon the detector 5 is comparatively increased which lowers the accuracy of analysis.\nAn interference-compensation type gas analyzer shown in FIG. 4 can improve an accuracy of analysis since the influences of the interfering component radiation are compensated by the subtractor 10.\nHowever, an interference-compensation detector 5b and a compensation-signal treatment circuit are excessively required, so that a problem occurs in an increase of cost."}
{"text": "The present invention relates to a process and an apparatus for the production of energy and a process and an apparatus for the production of methanol.\nNumerous processes and apparatus are already known for the production of energy and/or methanol from organic raw materials.\nSome of these processes and apparatus are based on the principle that an organic material is subjected to controlled oxidation wherein the energy carrier hydrogen is produced. The obtained hydrogen is supplied to a fuel cell and converted into current, according to the actual current requirements. Since storage of excess gaseous hydrogen is labourious and must especially satisfy stringent safety requirements when used in private households, hydrogen is normally catalytically converted together with carbon monoxide or carbon dioxide into the fluid energy carrier methanol in order to store its energy content, wherein the hydrogen can be re-extracted from the methanol and converted into current according to requirements by means of a reforming reaction.\nA process for the storage of energy present in the form of hydrogen is disclosed in DE 196 44 684 A1, wherein carbon dioxide, especially resulting from exhaust emissions, is converted after mixing with the hydrogen in a reactor into the energy carriers methane, methanol or ethanol.\nA process and an apparatus for the production of electrical energy from bio-raw materials is disclosed in DE 44 30 750 C2, wherein a hydrogen and carbon monoxide-containing crude fuel gas is produced from the bio-raw materials by means of partial oxidation with an oxygen-containing gasification medium in an oxidation reactor. From the fuel gas, the hydrogen proportion is placed into intermediate storage by reaction with a storage material, especially a so-called hydride store or with metal oxides, so that rapid removal of the hydrogen from the storage can be carried out according to requirements, and the hydrogen can be transported into a fuel cell module for conversion into current.\nA process and an apparatus for the production of hydrogen by means of the gasification of bio-raw materials is disclosed in DE 197 34 259 A1, wherein the bio-raw material is gasified into a hydrogen-containing crude gas in the presence of supplied water vapour in a vaporization reactor. The hydrogen proportion of the crude gas is separated from the residual gas at especially-high purity and can either be converted into current in a fuel cell, stored as methanol by synthesis with carbon dioxide or used for other purposes. In order to increase the degree of effectiveness in this process, the energy content of the low-hydrogen residual gas is also used in that it is combusted either in a closed cycle for the production of heat energy for the creation of water vapour or outside the cycle.\nA process is known from EP 0 257 018 A2 for thermal utilisation of waste and/or waste fuels, wherein the waste and/or waste fuels are supplied to a gasifying reactor which is directly heated with gas and/or oil and the gases extracted from the gasifying reactor are supplied to a combustion chamber for the production of energy, especially for the production of steam. A portion of the flue gases from the combustion chamber is fed back to the gasifying reactor under pressure.\nA process is known from DE 32 28 532 A1 for the carbonisation and gasification of solids which contain carbon wherein hot gases, amongst others, are used to heat the reactors, these gases being produced in a combustion chamber by combustion of a portion of the carbonisation gas produced in the carbonisation zone.\nIn xe2x80x9cEnergy Recovery From Waste: The Application of Gasification Technologiesxe2x80x9d, La Chimica e l\"\"Industria, 1996, No. 5, pp. 603-607, P. Pollesel gives an overview of the use of gasification processes to extract energy from waste, especially from bio-mass. Amongst other things, Pollesel mentions that the gases produced during gasification can be combusted in a gas turbine or a diesel engine to produce mechanical energy which can then be converted into electrical energy.\nConsequently, the object of the present invention is to provide a process and an apparatus for the production of energy and methanol respectively wherein optimised exploitation of the energy content of the organic starting material occurs in order to attain especially economical operation, and simultaneously allowing control of the proportions of released and stored energy respectively, according to requirements.\nA process according to claim 1 and an apparatus according to claim 9 or 11 is suggested as a solution for this object. The process for the production of energy, according to the invention, comprises the following steps:\n1. Combustion of fuel in an internal combustion engine to create mechanical energy and hot exhaust gas which contains carbon dioxide and water vapour,\n2. Reduction of the hot exhaust gas components carbon dioxide and water into a synthesis gas which contains carbon monoxide and hydrogen in an environment containing a supplied organic material,\n3. Supply of the synthesis gas produced in step 2 to the internal combustion engine.\nAccording to claim 7, methanol synthesis is inserted between steps 2 and 3 of the process according to the invention in order to produce methanol.\nThe steps of the process, according to the invention are set out for a recirculating operation, from which at different points energy in the form of electrical energy (current), heat energy and stored chemical energy (methanol) can be removed and/or added at various points, wherein the removal and supply can be regulated according to requirements, allowing optimised utilisation of the energy content of the organic starting material.\nTherefore in accordance with the invention, a fuel which is preferably fluid or gaseous and contains hydrocarbon is combusted on supply of air or oxygen in an internal combustion engine, creating mechanical energy, exhaust gas and heat energy. The exhaust gas released by the internal combustion engine, which contains carbon dioxide and water vapour, is supplied to a thermochemical reactor in which organic material such as biomass, coal, organic waste and such like is present. In the thermo-chemical reactor, carbon dioxide and water are reduced into carbon monoxide and hydrogen from the exhaust gas stream from the internal combustion engine, whilst the organic material is oxidised. The temperature required for the reaction is produced by the heat energy of the exhaust gas stream. Insofar as the temperature attainable by the inflowing exhaust gas is not sufficient to start or to maintain the reaction, a part of the hot exhaust gas stream from the internal combustion engine can be separated from the rest of the stream and used as an external heat supply for heating the thermochemical reactor and/or the temperature level required for the reaction can be increased by the supply of oxygen or air. Advantageously, the temperature of the hot exhaust gas is between approximately 900xc2x0 C. and approximately 1000xc2x0 C. In this way, the hot exhaust gas from the internal combustion engine is directly used as a valuable energy carrier whilst the mechanical energy produced by the internal combustion engine is available for other purposes.\nA synthesis gas which contains carbon monoxide and hydrogen is produced in the thermo-chemical reactor under suitable pressure and temperature conditions. The synthesis gas may contain other non-troublesome components such as, for example, carbon dioxide or methane, and undesired or troublesome solid, fluid and/or gaseous components such as, for example, suspended particles or sulphurous gases. The troublesome or undesired components can have a disadvantageous effect on further steps of the process, according to the invention, e.g. they can cause deactivation of a catalyst. Hence a filtration step is provided downstream for removal of the undesired or troublesome components.\nThe synthesis gas from the thermo-chemical reactor which contains carbon monoxide and hydrogen and which has been purified in the filtration step is mixed in a gas mixing unit with air oxygen or oxygen from an oxygen tank and is supplied to the internal combustion engine as a fuel. In this way, the supply of the thermo-chemical redox process with hot exhaust gas is ensured.\nInsofar as more synthesis gas is produced than is required for the operation of the internal combustion engine, as an alternative the synthesis gas, which contains carbon monoxide and hydrogen and which has been purified in the filtration process can be synthesised into the energy carrier methanol from carbon monoxide and hydrogen or carbon dioxide and hydrogen in an exothermic reaction according to a conventional process in a reactor which contains a suitable catalyst. Apart from heat energy, resultant reaction products are methanol and optionally water as well as non-converted synthesis gas. The reaction products are cooled in a downstream cooling unit. The precipitating methanol can be fed into a tank and stored. According to requirements, the stored methanol can be supplied to the internal combustion engine as a fuel, or used for other purposes e.g. as a fuel for motor vehicles.\nExcess, non-converted residual synthesis gas from the methanol synthesis is supplied to the internal combustion engine as described above.\nIn an embodiment of the Invention, the mechanical energy produced by the internal combustion engine is transformed into current by means of a generator. The current can either be available for the supply of a private household, stored in batteries or fed into the supply mains.\nIn a further embodiment of the invention, a fuel cell in a fuel cell unit is supplied via a gas mixer with oxygen and hydrogen in order to produce current, wherein the oxygen can be supplied from an oxygen tank or can be air oxygen and the hydrogen either originates from the synthesis gas purified by the filtration process or by reformation from the stored methanol. An especially-high degree of effectiveness for conversion of stored chemical energy into electrical energy can be attained with this type of current production. Furthermore, it is possible to adapt the production of current to an existing load by using the chemical energy of the stored methanol.\nIn an especially advantageous embodiment of the invention, heat resulting from the individual process steps, especially the radiated heat of the internal combustion engine, the generator, the methanol reactor and the cooling unit, can be supplied to a heat store where it is available, for example, for heating the thermo-chemical reactor or for heating rooms.\nIn a further especially advantageous embodiment of the invention, the material and energy streams can be controlled by means of a controller, e.g. a microprocessor control. This especially relates to the control of exhaust gas supply from the internal combustion engine into the thermo-chemical reactor, the control of the gas streams in the gas mixing unit, the monitoring and control of the parameters of the methanol synthesis, temperature regulation in various process steps, regulation of the energy streams of the heat storage and monitoring the level or regulation of the supply of organic material into the thermo-chemical reactor. In this way, optimum supply of the elements of the individual process steps with substrates and energy is attained at all times. Simultaneously, the energy content of the fuel is optimally used and adaptation of energy extraction and energy storage to the respective requirements is made possible.\nIn order to further solve the object which forms the basis of the invention, an apparatus is suggested for the production of energy and/or methanol, especially for implementation of the hereinbefore-described process. The apparatus, according to the invention, comprises an internal combustion engine, a thermo-chemical reactor disposed downstream of the internal combustion engine for reduction of the exhaust gas, which is produced in the internal combustion engine and which contains carbon dioxide and water steam, into a synthesis gas which contains carbon monoxide and hydrogen. An inlet is provided for the synthesis gas into the internal combustion engine so that energy may be produced. For the production of methanol, a methanol reactor is disposed downstream of the thermo-chemical reactor for Synthesis of methanol from the synthesis gas of the thermo-chemical reactor, wherein the synthesised residual gas and/or methanol from the methanol synthesis are fed into the internal combustion engine.\nIn an advantageous embodiment of the invention, a generator which can be driven by the internal combustion engine is provided for the production of current. Furthermore, a fuel cell apparatus disposed downstream of the thermo-chemical reactor can be provided for the production of current from methanol and/or at least part of the hydrogen produced in the thermo-chemical reactor.\nIn a further embodiment of the apparatus according to the invention, in order to store energy a heat storage is provided for storage of the heat generated with the apparatus. The heat storage serves especially for storage of the radiated heat from the exothermic process steps, for example from the internal combustion engine, the generator, the methanol synthesis and/or the cooling unit. The stored radiated heat is used in the exothermic process steps in the thermo-chemical reactor or for other purposes. Furthermore, a methanol tank is provided for storage of the synthesised methanol, wherein the stored methanol can be supplied to the internal combustion engine and/or the fuel cell apparatus if required.\nOther advantages and embodiments of the invention can be seen in the description and the accompanying drawing."}
{"text": "1. Field of the Invention\nThe present invention relates to an injection apparatus for melted metals used for injection molding nonferrous metals having a low melting point, such as zinc, magnesium, or alloy thereof, completely melted in liquid phase.\n2. Detailed Description of the Prior Art\nAttempts have been made to completely melt nonferrous metals having a low melting point so as to allow injection molding in liquid phase. Like in the case of injection molding of plastics, the molding method adopts a heating cylinder having inside an injecting screw, which is allowed to rotate and move along the axial direction. Granular metals supplied from the rear portion of the heating cylinder are heated and melted completely by shear heat and external heat while being transferred toward the fore end of the heating cylinder by means of rotation of the screw. After a quantity of the melted metals in liquid phase is metered in the fore portion of the heating cylinder, the metals are injected into a mold through the nozzle attached to the tip end of the heating cylinder by the forward movement of the screw.\nProblems occurring in case of adopting the foregoing injection molding for the metals are, for example, difficulty on the transfer of the material by means of rotation of the screw, the maintenance of the temperature of the melted metals in liquid phase, unstable metering, or the like.\nA melted plastic material has a high viscosity, and transfer of the melted plastic material by means of rotation of the screw is allowed mainly because a friction coefficient at the interface of the melted plastic material and the screw is smaller than a friction coefficient at the interface of the melted plastic material and the inner wall of the heating cylinder, and therefore, a difference in friction coefficient is produced between the two interfaces.\nIn contrast, the metal completely melted in liquid phase has such a low viscosity compared with the plastic material that a difference in friction coefficient is hardly produced between the above two interfaces. Hence, a transfer force such as the one produced with the melted plastic material by means of rotation of the screw is not readily produced.\nHowever, a transfer force is produced with the metals in solid state and in a high viscous region where the metals are in a semi-molten (liquid-solid) state during the melting process. Thus, the metals can be transferred by means of rotation of the screw up to that region. Nevertheless, as the metals are further melted, the viscosity thereof drops with an increasing ratio of the liquid phase, and the transfer force produced by the screw grooves between the adjacent screw flights decreases, thereby making it difficult to supply the melted metals in a stable manner to the fore end portion of the heating cylinder by means of rotation of the screw.\nBecause the melted plastic material has a high viscosity, it is stored in the fore end of the heating cylinder by means of rotation of the screw, while at the same time, a material pressure pushing the screw backward is produced as a reaction. By controlling the screw retraction caused by the material pressure, a constant quantity of the melted material can be metered each time.\nHowever, the metals in the low-viscous liquid phase cannot produce a pressure high enough to push the screw backward. Thus, the screw retraction by the material pressure hardly occurs, and if the metals are reserved in the fore end portion by means of rotation of the screw alone, a quantity thereof undesirably varies, thereby making it impossible to meter a constant quantity each time.\nIn addition, the metals have a far larger specific gravity compared with the plastics, and have a low viscosity and fluidity in liquid phase. For this reason, when allowed to stand by stopping rotation of the screw, the metals in liquid phase in the heating cylinder placed in a horizontal position leak into the semi-molten (liquid-solid) region in the rear portion through a clearance formed between the screw flights and the heating cylinder. Consequently, the metal material metered in the fore end portion causes a back flow onto the periphery of the fore portion of the screw through the opened ring valve, and the quantity thereof is undesirably reduced.\nThe liquid level in the fore end portion is lowered with the decreasing reserved quantity. For this reason, a gaseous phase (space) that makes the metering unstable is generated at the upper portion of the fore end portion. In addition, the leaked liquid phase material increases its viscosity in the semi-molten (liquid-solid) region as its temperature drops, or turns into solid depending on the heating condition in the semi-molten (liquid-solid) region, thereby forming weirs in the screw grooves. This poses a problem that the granular material supplied from the feeding opening provided behind the weir cannot be transferred readily by means of rotation of the screw.\nThe present invention is designed to solve the problems stated above in the injection molding of the metals in liquid phase. An object of the present invention is to provide a new injection apparatus which can easily and smoothly transfer the metals, melt them by the external heat, meter and degas by employing a reservoir to reserve metals in liquid phase for the injection screw, and a method for injection molding.\nIn order to achieve the above-mentioned object, the present invention according to the first aspect provides an injection apparatus for melted metals, comprising a heating cylinder having a fore end portion which communicates with a nozzle member and of which internal diameter is made smaller to serve as a metering chamber having a required length, and an injection screw installed within the heating cylinder to be movable and rotational, a tip end of the injection screw being formed in a plunger having a diameter which is almost the same as that of the metering chamber and can insert into the metering chamber while keeping a clearance for sliding, wherein a reservoir consisting of a portion is provided between the plunger and a feeding portion containing screw flight around the portion.\nMoreover, the present invention provides the injection apparatus for melted metals according to the foregoing aspect, wherein a projected portion for limiting the feeding of granular metals flowing from the feeding portion to the reservoir with metals in liquid phase and for preventing the metals in liquid phase reserved in the reservoir from flowing backward when the injection screw moves forward is provided on a boundary between said feeding portion and the reservoir.\nThe present invention further provides the injection apparatus for melted metals according to either of the foregoing aspects, wherein the screw flight of the feeding portion is provided in such a manner that screw groove of screw end is placed immediately below the feeding opening at the rearmost position of the screw in the heating cylinder, and that the screw end is placed in front of the feeding opening at the foremost position of the screw to close the feeding opening with the rear portion of the screw portion without screw flights, and to be capable of achieving transferring of the granular metals by the screw rotation at the rearmost position of the screw.\nThe present invention further provides the injection apparatus for melted metals according to the foregoing aspects, wherein the screw flight of the feeding portion is provided in such a manner that a screw groove of a screw end is placed immediately below the feeding opening at the foremost position of the screw in the heating cylinder, and that the screw end is placed behind the feeding opening at the rearmost position of the screw to be capable of achieving transferring of the granular metals by the screw rotation at the foremost position of the screw.\nMoreover, the present invention provides the injection apparatus for melted metals according to the first aspect, wherein the plunger is provided with a heat-resistant seal ring therearound, and a flow-through hole is formed therein from a ring groove for fitting the seal ring to a conical end of the plunger.\nThe present invention further provides the injection apparatus for melted metals according to any of the foregoing aspects, wherein the heating cylinder is installed with an inclination and positioning the feeding opening higher than the nozzle to allow the metals in liquid phase to flow down into the reservoir by its own weight.\nIn the construction stated above, a reservoir for the metals in liquid phase is provided between the plunger as a fore end portion and a feeding portion. By means of retracting the injection screw, the metal temporarily reserved in the reservoir is allowed to be reserved in the above-mentioned metering chamber. Thereby, the next feed of metals is completely melted and the temperature thereof is maintained while they are maintained in the reservoir even if the metals are melted by the external heat. As a result, the temperature of metals can be kept constant.\nSince a compressing portion to generate shear heat is unnecessary, the depth of the screw grooves between the screw flights can be made constant so as to feed the metals smoothly. Thereby the metals evenly contact the inner surface of the heating cylinder so that a fluctuation of temperature rarely happens. Since the most part of the metals melt into liquid phase while they reach to the projected portion on the boundary to the reservoir, and large granules which are incompletely melted are prevented from flowing into the reservoir by means of the projected portion, the metals in the reservoir are melted completely into the liquid phase and always ensured that they will be reserved into the metering chamber.\nFurthermore, in the construction stated above, while the screw moves forward and the feeding opening is being closed with the axis, the feeding of the metals will be automatically limited upon the start of injection. It prevents congestion of the metals in the screw grooves in the rear of the screw.\nThereby, a friction by rotation and sliding to the screw is decreased, which stabilizes melting and injecting of the metals to improve the quality of molded products.\nThe heating cylinder is inclined downward so as to reserve the melted metals in the reserving space surrounding the portion in the front portion of the heating cylinder. Therefore, even if the metals are in the liquid phase of a low viscosity, they will not flow backward so that the reserved amount will not fluctuate. In addition to it, since the rotation of the screw supplies the metals in liquid phase, in spite of injection molding the metals in liquid phase, a stable quality of molded metal products can be produced.\nThe nature, principle, and utility of the invention will become more apparent from the following detailed description when read in conjunction with the accompanying drawings in which like parts are designated by like reference numerals or characters."}
{"text": "This invention relates to a device for harvesting long agricultural products and an agricultural self-propelled unit for harvesting agricultural products comprising the device.\nThis invention is thus applied in particular in the agricultural sector, in particular in the harvesting of products performed after the cutting (or mowing).\nHarvesting devices, otherwise known as windrowers, are normally used for picking up from the ground grass, straw, hay (also cut by other machines) or for picking up pulses, as well as for picking up seed products, and, in any case, for picking up similar agricultural products, with a long shape.\nTypically, these devices are equipped with a suitable movement carriage and are connected to a tractor or the like which pulls them, like a trailer, along a direction of forward movement.\nThese devices are equipped with an arm, generally extending in a curved fashion from the rear portion to the front portion, to allow the connection with the rear hook of the tractor keeping the device oriented correctly, that is to say, with the pick-up means facing the tractor.\nOnly in some cases, for the small-sized devices (that is, equipped with a single pick-up unit), there is the possibility of connecting them to the front portion of the tractor, which, however, must be set up for this application, with considerable costs to be borne by the user.\nMoreover, it should be noted that this optional feature not may be acceptable for larger devices (that is, equipped with several pick-up units) since the visibility and driving problems would make it difficult to use the device.\nIn any case, from the above description it is evident that the prior art windrowers need a tractor for moving them, which represents a considerable cost for the user, especially if it is necessary to purchase a new pulling vehicle."}
{"text": "The present invention relates to connection members for components of a close-coupled pressurized system and, more particularly, a connector spool assembly provided with adjustment components to allow movement of the connector spool to facilitate separation and removal of system components.\nOne type of compression system is a compressor close-coupled to an electric motor driver, which provides for a compact design with significant benefits over traditional base-plate mounted compressor trains. A motor casing and a compressor casing comprise separate bodies requiring removal for service. One problem with component removal service activity is the cost and time required to disconnect process piping and instrumentation connected to each casing. Individual case removal is especially problematic for applications where the unit has compressor casings at each end of a double ended motor drive."}
{"text": "The present invention relates generally to imaging devices, and more particularly to a light guide for an array of detectors in an imaging device.\nIn certain types of imaging devices, such as positron emission tomography (PET) scanners, arrays of detector elements serve the function of detecting radiation emanating from the patient. In a PET scanner, for example, arrays of scintillator crystals detect gamma rays which are generated inside the patient. The gamma rays are produced when a positron emitted from a radiopharmaceutical injected into the patient collides with an electron causing an annihilation event. The scintillator crystals receive the gamma rays and generate photons in response to the gamma rays.\nOne of the challenges in designing a high resolution PET scanner relates to the space requirements of the electronics associated with the detector crystals, in particular the photomultiplier tubes (PMTs) which are situated behind the detector crystals. The function of the photomultiplier tubes is to receive photons produced by the scintillator crystals and to generate an analog signal with a magnitude representative of the number of photons received. The photomultiplier tubes typically cannot be diminished in size beyond a certain point, so that generally each photomultiplier tube is situated behind a number of smaller detector crystals. For example, a detector module in a PET scanner may comprise a 2×2 array of photomultiplier tubes situated behind a 6×6 array of scintillator crystals. In response to a scintillation event, each PMT produces an analog signal which is representative of the number of photons it has received. The relative magnitudes of the four PMT signals are then used to determine where the scintillation event took place and which crystal detected the event.\nIn determining the location of the scintillation event, it is generally advantageous to have a high degree of separation of the relative signal levels arising from each of the individual scintillation crystals in the detector array. Various arrangements have been proposed for increasing the spatial resolution of the detector crystals by controlling the light distribution within the detector array. For example the light distribution within the array of detector crystals can be controlled by applying various surface finishes having known light scattering and reflective properties to each crystal. These arrangements generally attempt to control the light distribution such that the proportion of light reaching each photomultiplier tube is relatively consistent and well defined for each event occurring at a particular detector crystal. In this way, the analog signals from the photomultiplier tubes may consistently determine which detector crystal produced the scintillation event.\nAs the demands for higher resolution in PET scanners continue to increase, one approach to achieving higher resolution is to increase the number of crystals in each detector array without increasing the size of the array. For example, a 6×6 array of detector crystals might be replaced with an 8×8 array. However, an increase in the number of crystals may introduce additional complexities and costs to the surface finishes and optical coupling which may be necessary for acceptable spatial resolution of the scintillation events. An increased number of smaller crystals may also introduce additional challenges with respect to light loss in the corner crystals and the tolerances for mechanical alignment of the array with respect to the photomultiplier tubes. The present invention provides an apparatus and method which can overcome these problems."}
{"text": "1. Field of the Invention\nThis invention relates to bistatic rotor tip synthetic aperture radars and more particularly to a relay link for transferring wide band information to or from a moving rotor tip and a second location wherein the relative motion provides a doppler shift to the wide band information requiring compensation.\n2. Description of the Prior Art\nRotor tip synthetic aperture radar (SAR) relates to radar mapping and surveillance systems in which transmit and/or receive apertures on a rotating radial arm utilize their motion to trace out long synthetic apertures in space from which radar signals can be sequentially transmitted and/or received so as to produce very fine angular resolution normal to the axis of rotation. Applications may be to helicopter rotors for navigation, obstacle avoidance, target detection, and other mapping functions and to fixed rotating radial arm towers for radar surveillance and ground testing and demonstration of synthetic aperture radar functions. Radar configurations may be monostatic with both transmission and reception from the rotating aperture or bistatic with only transmission or reception from the rotating aperture with the corresponding transmission or reception function performed by a fixed or non-rotating aperture.\nOne embodiment of a rotor tip SAR configuration may involve a radiating and/or receiving aperture at the tip of the radial arm, wave guide or coaxial transmission line down the arm, and microwave rotating joints through the arm rotation mechanism to the radar transmitter/receiver equipment at a fixed or non-rotating location. In applications to helicopters, especially, incorporation of rotating microwave joints in the transmission line down the blade or radial arm, may be constrained by considerations of flight performance and safety, and impose difficult and costly mechanical design problems. In addition, the long transmission line down the radial arm may impose a substantial loss in the signal level.\nIn communication systems where the transmitter is moving relative to the receiver, a doppler frequency shift according to equation 1 is observed. EQU f.sub.D =V f.sub.c /C (1)\nwhere f.sub.D =the doppler frequency shift, V represents the relative frequency between the transmitter and receiving stations, f.sub.c represents the frequency of the signal coupling the stations and C represents the velocity of light. Methods to compensate for the doppler frequency shift in communication systems have been described such as, for example, in U.S. Pat. No. 3,182,131, which issued on May 4, 1965 to R. R. Barnes. In U.S. Pat. No. 3,182,131, a pilot signal source having a predetermined frequency such as 62 kilohertz is superimposed on a carrier signal of 6 gigahertz along with information in the range from 562 kilohertz to 1298 kilohertz and transmitted to a receiving station. The receiving station has a local oscillator which is used to remove the carrier signal of 6 gigahertz. The pilot signal and information signals pass through respective filters. The pilot signal as received with the doppler frequency shift is multiplied by a predetermined amount to provide the doppler frequency shift of the information signals in the band. The doppler frequency shift is combined with a predetermined frequency from an oscillator in a mixer, which is presented to a second mixer, where the doppler frequency shift for the center of the predetermined channel is subtracted from the signals in the predetermined band, and thereby provides the information with a reduced frequency shift which is centered about the midpoint of the predetermined frequency band. In U.S. Pat. No. 3,182,131, the pilot signal travels with the information signal superimposed on a carrier signal.\nA doppler cancellation scheme is described in U.S. Pat. No. 3,325,736 wherein a signal is transmitted from a first station and received by a second station wherein it is processed to provide a signal with the frequency shift subtracted therefrom. The signal is then retransmitted from the second station to the first station wherein the doppler frequency shift added upon reception provides a signal where the first order of doppler frequency shift cancels. In U.S. Pat. No. 3,325,736, the transmitted and retransmitted frequencies may be sufficiently separated to permit continuous operation without feedback from the transmitting antenna to the receiving antenna of the same station. As shown in the drawing, each communication station must generate a predetermined frequency with respect to each other for proper operation.\nIn U.S. Pat. No. 3,450,842, a doppler frequency spread correction device is described which adds a doppler frequency to the base frequency.\nIt is therefore desirable to provide a relay link having one station in the rotor tip, and a second station mounted in a convenient structure, such as on the cab of a helicopter or on the supporting structure of the rotor for relaying radar related information from the rotor tip which may, for example, contain a receiver, to the cab of a helicopter which may, for example, contain a transmitter and signal processing equipment for processing a received signal.\nIt is further desirable to provide compensation for doppler frequency shifts arising from the relative motion of the first and second stations.\nIt is further desirable to transmit an auxiliary reference signal or pilot signal from the cab to the rotor tip and back to the cab where it may be used to cancel the one-way doppler frequency shift imposed on the radar data transferred from the rotor tip to the cab.\nIt is further desirable to phase lock the pilot signal with the radar signal being transmitted.\nIt is further desirable to offset in frequency the pilot signal being retransmitted at the rotor tip to the cab to facilitate simultaneous transmission and reception at both the cab and the rotor tip."}
{"text": "Advanced signal processing techniques are now being employed in the development of new television receivers. Such developments, termed ATV or advanced television techniques, include the concept of improved definition television, or IDTV, wherein the definition of the displayed image is enhanced by increasing the number of displayed lines over those present in the transmitted video signal. IDTV techniques become especially prominent in large screen sets where horizontal TV lines are discernable at normal viewing distances with present standards.\nHigh definition television techniques, on the other hand, involve methods that require changes in the transmitted signal, such as doubling the number of transmitted lines, and a complete redesign of the receiver. As such, HDTV standards impose an increase in the bandwidth of each channel beyond the limit set by the FCC (6 Mhz for the US NTSC standard). The stringent requirement by the FCC to remain within the bounds of the specified channel Bandwidth has caused delays and partial solutions to the evolution of a world HDTV standard. One of the recommendations for HDTV requested by many European countries, for example, is to move away from interlaced scanning to progressive scanning. Since IDTV works with the existing standards to provide an enhanced definition image, it represents an interim solution until an HDTV standard is agreed upon, and may also be used with HDTV to further enhance large displays. Even in the case of IDTV, however, very large intermediate storage arrays and very high-speed processing may be required to achieve the requisite level of signal processing, particularly if the goal is a system that operates in real time. An outstanding need, therefore, is a method and specialized hardware for achieving the requirements of IDTV."}
{"text": "1. Field of the Invention\nThis invention relates to steering wheel sensors that are mounted to rotating shafts that are joined by a torsion bar. In particular, there is an eccentricity compensator that prevents hysteresis between the two rotating shafts when their axes of rotation are co-axial or not co-axial.\n2. Description of the Related Art\nVarious devices and methods of dealing with the joining of two shafts that are rotating in a non-coaxial manner are known. These are called eccentricity compensators. One such device is a universal joint or U-joint. The U-joint is a well known device that typically is used to link between a transmission shaft and an axle shaft to allow the shafts angle to bend.\nUnfortunately, the prior art devices as they wear out, create what is called hysteresis. Hysteresis is a backlash or slop between the two rotating shafts. While, some hysteresis may be acceptable in a drive shaft application, in an application such as a steering wheel torque sensor this is unacceptable. A steering wheel torque sensor needs to have very precise position information of the relative rotational positions of the two shafts in order to correctly sense the amount of torque applied to the steering wheel. As the sensor wears, there cannot be excessive rotational movement between the two shafts leading to incorrect torque readings. Referring to FIG. 1, part of a prior art eccentricity compensator 10 is shown. Eccentricity compensator 10 is part of a steering wheel torque sensor that is described in U.S. patent application Ser. No. 09/837,075, filed Apr. 18, 2001 and titled, xe2x80x9cSteering Wheel Torque and Position Sensorxe2x80x9d.\nThe compensator 10 is shown in an assembled partial end view. The compensator 10 has a shaft 11 that is connected to a carrier or ring 14 by splines (not shown). Ring 14 has four pins 16 that extend upwardly. Another ring 12 is mounted adjacent to ring 14. Ring 12 has four slots 18. Pins 16 are located in slots 18. The rings 12, 14 and pins 16 are formed from injection molded plastic. Over a period of time during use, the eccentricity compensator parts will wear. The result is a gap 19 between pin 16 and a side wall 21. When the gaps 19 form in slots 18, the sensor components can rotate or have hysteresis when there is no actual torque in the steering wheel column. The compensator rotates about an axis or rotation 20. When the compensator is new rings 12 and 14 will rotate together in an original position as indicated by line 22. In compensator 10, there may be initial hysteresis due to a gap 19 due to manufacturing tolerances of the pin and slot. After the compensator 10 has had some wear, the positions of ring 12 and 14 will rotate relative to each other. This is indicated by dashed line 24. The rotational difference or error or hysteresis caused by the wear is indicated as a hysteresis angle 26. Angle 26 is very undesirable as it leads to incorrect torque readings.\nThe automotive industry has been focusing on electrical assist power steering for vehicles. The electrical assist power steering unit is an electrical motor attached to the steering linkage that operates when assist is required. A large amount of torque on the steering wheel occurs at low speed operation or during parking. The electrical assist power steering is generally not needed during high speed operation such as during highway driving. The major advantages of electrical assist power steering are first, that it only operates during the short time of turning and is inoperative the rest of the time and second that it is simpler to manufacture. In a hydraulic power steering system, the power steering pump is always being turned by the engine and represents an energy drain on the motor all the time even though steering is only performed during a small percentage of the total time a car is operated. An electrical assist power steering system requires sensing of torque applied to the steering wheel. The torque indicates how much force the operator is exerting to move the wheel. The output signal from a torque sensor is fed into a control unit which controls the electrical motor of the assist unit. When the torque sensed is high, the assist applied to the steering linkage will be high. When the torque sensed is low, the assist applied to the steering linkage will be low.\nIn general, a sensor that measures the relative displacement between two rotating shafts has useful applications in the areas of industrial machinery, aerospace, electrical power generation and transportation.\nThere is a current unmet need for a device that prevents hysteresis between two non-coaxial rotating shafts. Additionally, there is a current unmet need for a eccentricity compensator to prevent hysteresis in a steering wheel column torque sensor.\nIt is a feature of the invention to provide a eccentricity compensator that prevents wear induced rotational displacement that is mounted between two rotating shafts that have non-coaxial axes of rotation.\nYet, another feature of the invention is to provide a eccentricity compensator for preventing hysteresis between a first and a second rotating shaft that are joined by a torsion bar. The eccentricity compensator includes a first ring attached to the first shaft and a second ring located adjacent the first ring. A third ring is attached to the second shaft and the third ring is located adjacent the second ring. At least one pin-slot pair is located between any two of the first, second or third rings. A spring is mounted adjacent the pin-slot pair. The spring biases the pin-slot pair such that hysteresis is prevented. The springs and slots are arranged in such a way that wear of the pins in the slots does not result in drift.\nThe invention resides not in any one of these features per se, but rather in the particular combination of all of them herein disclosed and claimed. Those skilled in the art will appreciate that the conception, upon which this disclosure is based, may readily be utilized as a basis for the designing of other structures, methods and systems for carrying out the several purposes of the present invention."}
{"text": "The present invention relates to a process for detecting live microbiological contaminants in a food product sample.\nFood products are particularly susceptible to contamination with microbiological products, in particular with bacteria, on account of the powerful effect which a contamination has on the health of the persons ingesting the food. In this sphere, live, and therefore active, microorganisms are particularly formidable since they are then able to propagate in the body and transmit severe diseases. There is therefore a definite need to develop a process which enables these live microorganisms to be detected in food products in a manner which is both precise and satisfactory.\nA large number of methods exist for detecting bacteria, in particular, in samples. One which may be mentioned, by way of example, is that of culturing the sample in order to increase the number of bacterial cells present until colonies, which can be counted, are observed. Where appropriate, the culture can be followed by an additional step which enables the particular type of bacteria contained in the sample to be identified. These bacteriological methods require a great deal of time and skill on the part of the individuals who carry them out. Thus, it is generally necessary, for example, to incubate for from 24 to 48 hours before being certain of obtaining a positive or negative result.\nOther methods have also been envisaged for eliminating the drawbacks of the conventional bacteriological method. Thus, microscopy is frequently used for detecting bacteria in clinical samples. Usually, it is necessary to stain the sample in order to increase the detection limits, a procedure which on the one hand represents a laborious method and, on the other hand, is unsuitable for food samples (1). Some immunological methods have also been successfully developed for detecting certain species which possess surface antigens which can be recognized by specific antibodies. However, such an approach cannot be used for qualitatively determining bacteria in a sample which may contain a large variety of bacterial species which do not possess a common surface antigen.\nThe European Patent Application EP-A-0 133671 describes a method for determining the presence of bacteria in samples, in particular in media, such as body fluids, which are suitable for diagnostic purposes, which employs nucleic acid hybridization techniques. According to this method, the sample to be tested is first of all subjected to denaturing conditions so as to render the nucleic acids of the bacteria, which are present in the sample, single-stranded. The resulting single-stranded nucleic acids are brought into contact with a polynucleotide probe which possesses a sequence which is homologous with at least one sequence which is common to all the bacterial species which are present in the sample. The probe and the denatured nucleic acids are brought into contact so as to effect a hybridization between the probe and the respective sequences of the bacteria. More particularly, the probe employed comprises at least a part of one of the strands of a gene which codes for the synthesis of a nucleic acid or a protein which is involved in the mechanism by which the proteins are synthesized. Those genes of this type which are cited are, in particular, genes which encode transfer RNAs, ribosomal RNAs, or initiation, elongation or translation termination factors. One of the features of this method is that it is not based on expression but rather on the presence of hybridizable nucleic acids, for example the RNA or the genomic DNA or the extra-chromosomal nucleic acids of the bacteria which are present. The result, which the Applicant views as being an advantage, of this is that the samples do not have that it is not based on expression but rather on the presence of hybridizable nucleic acids, for example the RNA or the genomic DNA or the extra-chromosomal nucleic acids of the bacteria which are present. The result, which the Applicant views as being an advantage, of this is that the samples do not have to be treated so as to guarantee the viability of the bacteria which are present.\nDocument WO 92/03455 describes compositions and processes for treating and diagnosing infections with Candida, in particular nucleotides which are able to hybridize specifically with a part of the Candida MRNA, in particular the mRNA encoding elongation factors 1 and 2 (TEF1 and TEF2). This document is only directed towards therapeutic or diagnostic applications, either for detecting the presence of Candida in a patient or for inhibiting the activity of this bacterium by blocking the expression of essential proteins.\nThe document Berg et al., xe2x80x9cMOL CELL. PROBES, Vol. 10, February 1996, pp 7-14xe2x80x9d describes a system for specifically detecting the DNA of microplasmas using the PCR technique. Although it is indicated, on page 12 of this document, that the target sequence for the PCR amplification is the tuf gene, which encodes the elongation factor Tu, this document is only directed towards detecting DNA, and not mRNA on the one hand, and, on the other hand, the method is only a method for detecting bacteria in order to establish a diagnosis in a patient. This document does not envisage any application in the sphere of the invention, in which the specific problem is that of detecting living bacteria.\nThus, there is no known method which makes it possible to detect live microbiological contaminants in a food product and which at the same time discriminates between the live microorganisms and the dead microorganisms and which, moreover, does not pose any problems relating to public health.\nThe inventors have now discovered that it was possible to detect live microbiological contaminants in a food product sample by detecting, in this sample, the resence of messenger RNA (mRNA) which codes for the synthesis of a protein which is involved in the mechanism by which the proteins of the said contaminants are synthesized.\nDifferent families of microorganisms can be detected in accordance with the invention. Procaryotes, in particular bacteria, unicellular eucaryotes, in particular yeasts, and multicellular eucaryotes, in particular fungi, may be mentioned in a nonlimiting manner. Different species can be identified within these families. Thus, for example, Escherichia, Salmonella and Mycobacterium in the case of bacteria; Saccharomyces and Candida in the case of yeasts; Mucor, Neurospora and Trichoderma in the case of fungi.\nThe synthesis of proteins by microorganisms comprises steps of transcription and translation. Within the context of translation, nucleic acids and proteins exist which are involved in each of the three basic steps of protein synthesis, i.e. initiation, elongation and termination. invention in order to detect live microbiological contaminants belonging to different species.\nThus, according to one preferred embodiment, the invention relates to a process for non-specifically detecting live contaminants belonging to different species of a family of microorganisms, according to which the MRNA detected is an mRNA which codes for the synthesis of a protein whose primary structure is at least partially conserved between different species.\nOn the contrary, if a greater specificity of detection between different microorganism species is required, it is possible, according to the invention, to detect different mRNAs which respectively code for the synthesis of a protein whose primary structure is not conserved between different species.\nOf all the factors which are involved in protein synthesis, one example which is particularly preferred consists of the elongation factors. Those of these factors which may be mentioned are the EF-1, EF-2, EF-G and EF-TU factors, in the case of bacteria (2), or else the EF-1xcex1 factor (3) in the case of yeasts and fungi. These factors play a fundamental role in protein synthesis in that they determine the length of time during which an aminoacyl tRNA remains associated with the ribosome and with the forming polypeptide chain, which function is referred to by the expression xe2x80x9ckinetic proofreadingxe2x80x9d (4).\nThere are various reasons why it is particularly advantageous to look for the presence of messenger RNA which encodes an elongation factor. First of all, this gene represents a very suitable marker of the viability of the cells since inactivation of this gene is a lethal event both in procaryotes and eucaryotes (5, 6). Furthermore, the gene which encodes an elongation factor encodes a protein which belongs to those proteins which are most widely expressed in procaryotes and eucaryotes (7, 8), a fact which makes it possible to substantially decrease the cell detection level. Finally, as pointed out above, it is possible to modulate the specificity of the detection insofar as this function is conserved in procaryotes and eucaryotes and the primary structure of this type of gene is very similar (9). It is thus possible to implement means of detection which make it possible to distinguish between procaryotes and eucaryotes or, on the contrary, to non-specifically detect bacteria, yeasts and/or fungi at one and the same time.\nFurthermore, mRNA encoding elongation factors has a very short half-life (10, 11). Its presence therefore reveals the presence of cells which were still alive approximately ten minutes before the mRNA was detected.\nAn mRNA encoding an elongation factor is therefore detected in accordance with a preferred embodiment. In this case, a live microorganism cell within the meaning of the invention is a cell which is able to produce the mRNA corresponding to an elongation factor.\nThe RT-PCR (polymerase chain reaction combined with reverse transcription) is a method of choice for detecting the presence of messenger RNA according to the invention. This technique consists in carrying out a PCR on an RNA which has previously been transcribed into complementary DNA in the presence of reverse transcriptase and a primer. After the RT stage, the proper PCR stage is carried out under standard conditions in the presence of the DNA to be amplified, two oligonucleotide primers which flank the region to be amplified and four deoxynucleotide triphosphates (DATP, dCTP, dGTP and dTTP), in large molar excess, and the enzyme Taq polymerase.\nNaturally, the choice of the primers is a basic requirement, since it makes it possible to target the mRNA which it is desired to detect. The oligonucleotide primers are prepared such that they are specific for the coding region of a gene which encodes an elongation factor. The known elongation factors which are preferably selected are the EF-TU factor in the case of bacteria and the EF-1xcex1 factor in the case of yeasts and fungi. It is for this reason that the primers B1/B2 (5xe2x80x2 CGCTGGAAGGCGACGMRRAG 3xe2x80x2 (SEQ ID NO:1)/5xe2x80x2 CGGAAGTAGAACTGCGGACGGTAG 3xe2x80x2 (SEQ ID NO:2) were, for example, prepared, which primers are specific for a fragment of the bacterial EF-TU elongation factor which is found, in particular, in Salmonella typhimurium, Mycobacterium tuberculosis, Mycobacterium leprae, Escherichia coli, Brevibacterium linens and Streptomyces ramocissimus. \nIn the case of yeasts, the primers L1/L2 (5xe2x80x2 TCCATGGTACAAGGGTTGGGAA 3xe2x80x2 (SEQ ID NO:3)/5xe2x80x2 GCGAATCTACCTAATGGTGGGT 3xe2x80x2 (SEQ ID NO:4) were prepared, which primers are specific for a fragment of the yeast EF-1xcex1 elongation factor, which is found both in Saccharomyces cerevisiae and Candida albicans. \nFinally, an example of a nucleotide primer which can be used for detecting messenger RNA which is specific for fungi consists of the pair M1/M2 (5xe2x80x2 GCTGGTATCTCCAAGGATGG 3xe2x80x2 (SEQ ID NO:5)/5xe2x80x2-CGACGGACTTGACTTCRGTGG 3xe2x80x2 (SEQ ID NO:6). These primers are more particularly specific for a fragment of the fungal EF-1xcex1 elongation factor which is found in Mucor racemosus, Neurospora crassa, Trichoderma reesei, Absidia glauca, Aureobasidium pullulans, Histoplasma capsulatum and Puccinia graminis. \nThe primers which can be used in accordance with the invention were prepared after comparing the elongation factor-encoding regions of different microorganism species (Lasergene software, Dnastar, Madison, Wis., USA).\nIn a general manner, the process according to the invention can be characterized by the following steps taken together:\na) a sample of the food product to be tested is withdrawn;\nb) the cells are lysed;\nc) reverse transcription is carried out;\nd) amplification cycles are carried out using oligonucleotide primers which are specific for the coding region of a gene which encodes an elongation factor;\ne) the amplification products are separated;\nf) the amplification products are visualized.\ng) the products as visualized in step f) are compared with the amplification products which are obtained from pure mRNA.\nSince the invention is essentially directed towards detecting live microbiological contaminants, it is particularly advantageous to as far as possible remove any additional contamination of the medium with the DNA which is present in the sample. For this reason, according to a preferred embodiment of the invention, an additional step (bxe2x80x2), which is intended to remove the DNA which is present in the sample, is added after step b). For this, DNase I, which does not contain any RNase, can simply be added to the reaction medium. If it is desired to check for the absence of false positives which are linked to the presence of DNA, one possibility consists in carrying out a PCR reaction on the same samples as those used for the RT-PCR."}
{"text": "1. Field of the Invention\nAspects of the present invention relate to a sealing apparatus and a method of manufacturing a flat display device using the sealing apparatus.\n2. Description of the Related Art\nDisplay devices are increasingly being replaced by slim, portable flat panel display devices. Electroluminescent display devices, among the flat panel display devices, are self luminescent display devices, and have advantages of a wide angle of view, a high contrast level, and a fast response speed. Therefore, the electroluminescent display devices are being pursued as a display device of the next generation.\nOrganic light-emitting display devices having a light-emitting layer composed of an organic material have characteristics of brightness, a lower driving voltage, and a higher response speed and are capable of displaying more colors than inorganic light-emitting display devices. Generally, an organic light-emitting display device has a structure in which at least one organic layer including an emission layer is interposed between a pair of electrodes (that is, a first electrode and a second electrode). The first electrode is formed on a substrate, and functions as an anode that injects holes. An organic layer is formed on the first electrode. The second electrode functions as a cathode that injects electrons, and is formed on the organic layer so as to face the first electrode.\nWhen moisture or oxygen flows into elements of such an organic light-emitting display device from the surroundings, there are problems due to oxidation of an electrode material. Examples of such problems are exfoliation of an electrode surface or the like. These problems shorten the life of the elements, lower a light-emitting efficiency, and deteriorate emitted light colors. Therefore, in the manufacture of the organic light-emitting display device, a sealing process is usually performed so that the elements can be isolated from outside exposure and moisture infiltration into the elements. In the sealing process, usually a polymer organic material (such as polyethylene terephthalate (PET)) that is a thermoplastic polyester, is laminated on the second electrode or a cover or a cap is formed on the second electrode. Nitrogen gas is filled in the cover or the cap, and the rim of the cover or the cap is capsule-sealed by a sealant such as epoxy. The cover or the cap may be made of a metallic material or a glass material that is moisture absorbent.\nHowever, such methods can allow for some element damaging materials such as moisture, oxygen, or the like to infiltrate into the device from the outside. In addition, the methods are not well suited when applied to top-emitted organic light-emitting display device having elements that are specifically weak to moisture, and a process to realize the methods is complicated. In order to solve the above mentioned problems, a capsule sealing method is used to enhance adhesion between an element substrate and a cap using a frit as a sealant.\nThe method of capsule-sealing an organic light-emitting display device includes coating a frit on a glass substrate which perfectly seals the element substrate and the cap together. Thus, the organic light-emitting display device can be protected from being infiltrated with the element damaging material more effectively than an organic light-emitting display device using epoxy as a sealant.\nIn the method of capsule-sealing the organic light-emitting display device with the frit, after the frit is coated on each sealing portion of the organic light-emitting display device, a laser irradiating apparatus is moved to the sealing portions thereof and irradiates a laser beam to the sealing portions to harden the frit for sealing."}
{"text": "The present disclosure relates to a display unit displaying an image, and an image processing unit for use in such a display unit, and a display method.\nRecently, a cathode ray tube (CRT) display unit has been actively replaced with a liquid crystal display unit or an organic electro-luminescence (EL) display unit. The liquid crystal display unit and the organic electro-luminescence display unit are each being a mainstream display unit due to low power consumption and a flat configuration thereof.\nDisplay units are in general desired to have high image quality. Image quality is determined by various factors including contrast. Increase of peak luminance may be a technique for improving contrast. Specifically, reduction of a black level is limited by reflection of outside light, etc. Hence, in the above technique, peak luminance is increased (extended) to improve contrast. For example, Japanese Unexamined. Patent Application Publication No. 2008-158401 (JP-A-2008-158401) discloses a display unit, in which an increasing level (extending level) of peak luminance and gamma characteristics are each varied depending on an average of image signals to achieve improvement in image quality and reduction in power consumption.\nIn some display units, each pixel is configured of four sub-pixels. For example, Japanese Unexamined Patent Application Publication No. 2010-33009 discloses a display unit, in which each pixel is configured of sub-pixels of red, green, blue, and white to improve luminance or reduce power consumption, for example."}
{"text": "This invention relates generally to plural-substituted pyridine derivatives and, particularly, to processes for preparing and isolating 5-trifluoromethyl-2-pyridone and for using same.\nAlthough a large number of functionally substituted pyridine compounds are known and capable of synthesis, certain patterns of disubstitution on the pyridine ring are difficult to obtain by any convenient and commercially viable means. Pyridines having functional substituents in the 2- and 5- positions of the ring are often valuable derivatives, but fall within this category. For example, hydroxyl, cyano, carboxy, chloro and other groups are difficult to introduce into these positions on the pyridine ring.\nIt has long been known that certain non-nitrogen-containing heterocycles can assist in this regard because of their ability to be converted into pyridine bases. Pyrones, pyrilium salts and furans are examples of such transformations. As early as 1884 in Ber., 17, 2384 (1884), Von Pechmann et al. described the conversion of 5-carboxy-2-pyrone to 5-carboxy-2-hydroxypyridine upon treatment with ammonia in the presence of a caustic material such as sodium hydroxide. ##STR2## Although this reaction is potentially useful in some commercial applications, the number of reported uses of the method is small.\nThis 5-carboxy-2-hydroxypyridine, also known as 6-hydroxynicotinic acid and 5-carboxy-2-pyridone, is an example of a disubstituted pyridine of some commercial value. Besides this Von Pechmann et al. method, it has been prepared directly from a pyridine in two known instances. First, synthesis has been achieved by direct carboxylation of 2-hydroxypyridine, as depicted below. ##STR3## As can be seen, this procedure is rather lengthy and has proven not a viable commercial method. Second, 2-hydroxy-5-carboxypyridine has been prepared directly from 5-carboxypyridine, also known as niacin, through known techniques.\nTrifluoromethyl-substituted pyridines have also proven to be valuable derivatives of pyridine bases, although difficult to obtain. Synthesis of these compounds has been accomplished in only a limited number of cases, most of which have involved the conversion of a pyridine carboxylic acid to a trifluoromethyl pyridine utilizing sulfur tetrafluoride.\nAccordingly, it is generally known that alkyl and aromatic carboxylic acids react with sulfur tetrafluoride in the presence of hydrogen fluoride to give trifluoromethyl derivatives. ##STR4## This reaction has been applied to amino acids. Kobayashi et al.: Chem. Pharm. Bull. 15, 1896 (1967), M. S. Raasch: J. Org. Chem., 27, 1406 (1962). It has also been applied in a few reports to simple pyridine carboxylic acids such as niacin and 3,5-dicarboxypyridine. Kobayashi et al. and Raasch, id.\nIt is likewise generally known that some esters and anhydrides of these carboxylic acids react with sulfur tetrafluoride to give the corresponding fluorinated ethers. In the case of compounds such as ethyl acetate and dichloromaleic anhydride, the double-bonded oxygen groups are simply replaced during the fluorination reaction. W. R. Hasek, W. C. Smith, and V. A. Engelhardt, J. Amer. Chem. Soc., 82, 543 (1960). ##STR5##\nThe presence of a hydroxy group leads to undesirable reactions involving this additional oxygen function in the trifluoromethylation reaction. As used herein, \"trifluoromethylation\" refers to the conversion of a precursor material to a material containing a trifluoromethyl radical by the addition or substitution of fluorine to the precursor material. As background for this statement, 2- and 4-hydroxypyridines have been shown to physically exist as a mixture of tautomeric forms, appearing both as the hydroxy and the amide derivatives. R. Elderfield, Heterocyclic Compounds, 1, 435-440 (1950). ##STR6## For this reason, these hydroxypyridines undergo reactions typical of both phenols and amides as also reported in the Elderfield reference. ##STR7## Similar behavior has been reported in substituted hydroxypyridines such as the 5-carboxy-2-hydroxypyridines discussed above. Klingsberg, Pyridine and Its Derivatives, Part Three, p. 646 (1962).\nIt is known that hydroxy groups give rise to fluoro groups by standard substitution upon treatment with sulfur tetrafluoride. Boswell et al., Org. Reactions, 21 p. 12 (1973). It is also known that amides react with sulfur tetrafluoride to give a variety of products. Hasel, Smith and Engelhart, J. Amer. Chem. Soc., 82, 543, (1966) For example, if the amide contains at least one nitrogen-hydrogen bond, cleavage at the nitrogen-carbon bond is reported to occur. ##STR8## This cleavage is believed to be caused by trace amounts of hydrogen fluoride produced during the reaction. Hasek, Smith & Englehart, J. Amer. Chem. Soc., 82, 543 (1966).\nTherefore, disubstituted 2- and 5-pyridine derivatives are often difficult to obtain. This statement is particularly true, as taught by the art, with the 2-hydroxy and 5-trifluoromethyl substituents. Nevertheless, 5-trifluoromethyl-2-pyridone is now proving to be a desirable and valuable commercial compound based on both proven and anticipated uses as shown in the art. It appears from the prior art that 5-trifluoromethyl-2-pyridone would be useful as a catalyst in nucleophilic aromatic substitution, formation of amides, and hydrolysis of esters. It also appears from the prior art that 5-trifluoromethyl-2-pyridone would be useful as an antioxidant. Still further, this compound is proving to be a valuable intermediate in the synthesis of herbicides, pharmaceuticals, germicides and the like. See U.S. Pat. No. 4,038,396 to Shen et al. This is due at least in part to the ready substitution for the 2-hydroxy group on the ring and the lower toxicity caused by the 5-trifluoromethyl substituent.\nThe only reported synthesis of 5-trifluoromethyl-2-pyridone in the art attempts to avoid the prior art problems discussed above. This synthesis, reported in U.S. Pat. No. 4,038,396 issued to Shen et al. on July 26, 1977, teaches a three-step method in Example 111 (beginning at column 23, line 62) for preparing 5-trifluoromethyl-2-pyridone (the tautomeric form of 5-trifluoromethyl-2-hydroxypyridine) from 6-hydroxynicotinic acid. This method includes steps for specifically circumventing the expected interference of the tautomeric hydroxy and amide forms of the initial compound. This is accomplished by converting the 2-substituent to a chloro group in an attempt to protect the 2- position during the trifluoromethylation reaction.\nIn particular, Shen et al. first teach converting this 6-hydroxynicotinic acid to a less-reactive 6-chloronicotinic acid by a reaction involving liquid phosphorous oxychloride added jointly with solid phosphorous pentachloride and then recaptured in water. The 6-chloro derivative is then subjected to trifluoromethylation using sulfur tetrafluoride in the presence of hydrogen fluoride, and the 2-chloro-5-trifluoromethylpyridine is reconverted back to the hydroxy or amide derivative through a complex treatment under nitrogen with silver acetate in the presence of acetic acid.\nThis three-step Shen et al. procedure is lengthy and complex, and thus of minimal commercial valuable. It further emphasizes the need for the development of a viable, more efficient process for preparing the valuable 5-trifluoromethyl-2-pyridone. ##STR9##"}
{"text": "1. Field of the Invention\nThe present invention relates to a recording medium, a substrate for a recording medium, production processes thereof, and an image forming process using the recording medium. In particular, the present invention relates to a recording medium, which can provide a clear or bright and high-quality recorded image in a surface coated part region thereof and can prevent the occurrence of a phenomenon called cockling in which a printed surface is waved by an aqueous ink, and a production process of the recording medium.\n2. Related Background Art\nAn ink-jet recording system often used in recent years is a system in which fine droplets of an ink are flown by any one of various working principles to apply them to a recording medium such as paper, thereby making a record of images, characters and/or the like. Recording apparatus, to which this recording system is applied, are quickly spread as recording apparatus for various images in various applications including information instruments because they have features that recording can be conducted at high speed and with a low noise, multi-color images can be formed with ease, printing patterns are very flexible, and neither development nor fixing is unnecessary. Further, they begin to be widely applied to a field of recording of full-color images because images formed by a multi-color ink-jet system are comparable in quality with multi-color prints by a plate making system and photoprints by a color photographic system, and such images can be obtained at lower cost than the usual multi-color prints and photoprints when the number of copies is small.\nWith the enlarged utilization of the ink-jet recording system, further improvements in recording properties such as speeding up and high definition of recording, and full-coloring of images are required, so that recording apparatus and recording methods are improved. In order to meet such requirements, a wide variety of recording media have heretofore been proposed. For example, there have been proposed paper for ink-jet recording, in which a coating layer having good ink absorbency is provided on a surface of a substrate (for example, Japanese Patent Application Laid-Open No. S55-5830), and the use of amorphous silica as a pigment in an ink-receiving layer laminated on a substrate for recording medium (for example, Japanese Patent Application Laid-Open No. S55-51583).\nWith the diversification of uses, it has also been required to reduce the occurrence of curling or cockling of printed articles for the purpose of improving the quality of recorded images. These phenomena are both considered to be caused due to the occurrence of expansion or shrinkage and distortion on a recording medium by absorption of an ink. In the present invention, the cockling means a phenomenon that a printed surface of a recording medium is made irregular or waved. As means for avoiding this cockling phenomenon, there have heretofore been proposed the following methods.\n(1) Japanese Patent Application Laid-Open Nos. H3-38376, H3-199081, H7-276786 and H8-300809 describe recording media using paper having an underwater elongation and a wetted elongation within respective specified ranges. However, since the technical ideas described in these documents are based on the premise that water is evenly given to the whole of a recording medium, they cannot cope with a case where states applied with a liquid differ with portions.\n(2) Japanese Patent Application Laid-Open No. H10-46498 discloses a crosslinking treatment in which a water-proofing agent, a polymer, a size and the like are used to form a bound structure between fibers, and also discloses to the effect that the degree of floating after 10 seconds from printing is controlled to 1 mm or small. Japanese Patent Application Laid-Open No. 2002-201597 has proposed a recording medium in which cellulose fiber is shrunk by a mercerization treatment that a treatment with an alkali is conducted on the whole surface, and discloses to the effect that friction with an ink-jet recording head is avoided. Incidentally, these proposals are both those for recording media on which no ink-receiving layer is provided.\n(3) The constitution that an ink-receptive layer containing a water-repellent component in Japanese Patent Application Laid-Open No. 2000-158805 and a void layer formed of a thermoplastic resin such as polyurethane in Japanese Patent Application Laid-Open No. 2002-154268 are respectively provided as intermediate layers for barrier preventing penetration of inks between an ink-receiving layer and a substrate is described. Since these intermediate layers both act as a barrier preventing penetration of inks, the quantity of inks absorbed is reduced, and an ink-absorbing speed is lowered when the quantity of inks printed is great because the inks printed do not penetrate into the substrate, so that ink overflowing and/or bleeding may be caused in some cases.\n(4) Proposals for the solution, which are different from the methods in the above-described publicly known documents, include the following proposals. Namely, the proposals comprise providing an additional structure on a recording medium. A recording medium, in which ink-receptive layers are provided on both surfaces of a substrate, a recording medium, in which a back coat layer is provided on a surface opposite to an ink-receiving layer, and a recording medium, in which substrates are laminated on each other into a two-layer structure, are described in Japanese Patent Application Laid-Open Nos. H2-270588, 2001-253160 and 2002-2092, respectively.\n(5) On the other hand, Japanese Patent Application Laid-Open No. 2002-211121 discloses a recording medium, in which an aqueous solution containing a cationic resin and an alumina hydrate is coated on an ink-receiving surface of a single-layer fibrous structure composed mainly of a fibrous material containing no filler and making no use of a size (non-sized). It is described that according to the structure disclosed in this document, the cationic resin and alumina hydrate can be caused to exist on the surface of the fibrous material, thereby surely trapping an anionic colorant, so that an excellent image can be formed without causing cockling and very great curling upon formation of a 100% solid-printed image.\n(6) Japanese Patent Application Laid-Open No. 2001-246840 discloses a recording medium, in which an ink-receiving layer having an inorganic pigment and a binder is formed in a coating weight of 1 to 10 g/m2 on a base material composed mainly of pulp fiber. In Comparative Example 2, it is described that when an ink-receiving layer containing no binder was formed on a base material subjected to a sizing treatment, an alumina hydrate that is an inorganic pigment entered pulp fiber, and the surface of the pulp was scarcely coated with the alumina hydrate.\nThe present inventors have carried out an investigation on various kinds of the recording media proposed in the prior art documents mentioned above and found, on all the recording media, a phenomenon that new cockling or curling is caused when printing is conducted in an ink quantity exceeding 100% in particular. When the state thereof has been analyzed, it has been found that the number of portions undergoing cockling substantially increases, so that the cockling is conspicuous. The present inventors have also found that when a quantity of an ink applied to a recording medium is increased to 2 times or 3 times to form an image, the ink-absorbing capacity of the recording medium itself is lowered, so that ink overflowing and/or bleeding may be caused in some cases to fail to achieve good image quality.\nThe present inventors have paid attention to the fibrous materials of the substrates to carry out research and investigation. As a result, the following facts have been confirmed. Since the alumina hydrate and cationic resin are coated on the fibrous material by on-machine coating in Japanese Patent Application Laid-Open No. 2002-211121, the alumina hydrate is limited by the application of the cationic resin to the surface of the fibrous material and partially scattered. It has been confirmed that the cockling-inhibiting effect by the alumina hydrate is not sufficiently brought about on the recording medium disclosed in this document as described below.\nIn Japanese Patent Application Laid-Open No. 2001-246840, the substrate containing pulp and a filler and size-pressed is used. Accordingly, when the alumina hydrate is applied without using a binder like Comparative Example 2 of this document, the alumina hydrate cannot be applied to the surface of the substrate fiber, but only fills in voids formed by the pulp. It has been thus confirmed that this constitution does not bring about the cockling-inhibiting effect as described below.\nAs described above, it has been confirmed that when images are formed by printers for conducting high-speed printing in recent years, or the like, even the various kinds of recording media proposed in the prior art documents are not always satisfied from the viewpoints of image quality, curling, cockling, conveyability and the like.\nIt is a principal object of the present invention to solve the novel problems on the basis of such new findings.\nThe present inventors have sought a phenomenon by deformation of fiber, such as swelling and elongation, that is a cockling producing mechanism and considered that it is caused by excessive absorption of water by the fiber and a high degree of freedom of displacement within an allowable space. Accordingly, the present inventors have sought means that a water-holding capacity of the fiber itself can be diffused so as to be optimized, and at the same time the degree of freedom of displacement can also be controlled, thus leading to completion of the present invention.\nIt is thus a first object of the present invention to provide a recording medium having an ink-receiving layer that can solve the above-described new problem caused by conducting printing in an ink quantity exceeding 100%, permits forming an image high in density and bright in color tone and can settle the cause of the occurrence of new curling or cockling, and a production process of the recording medium.\nA second object of the present invention is to provide a substrate (including a case where the substrate itself functions as an ink-receiving layer) for a recording medium for preventing ink overflowing even in an ink quantity exceeding 100%, permitting forming an image high in density and bright in color tone and settling the cause of the occurrence of new curling or cockling, and a production process of the substrate for a recording medium.\nThe above objects can be achieved by the present invention described below."}
{"text": "This invention relates to soldering.\nIn one method of soldering, printed circuit boards (PCBs) which are populated with components pass, one at a time, through a three step process: flux is applied to electrical connection points on both the PCB and the components; the PCB and the components are preheated; and the electrical connection points are brought in contact with molten solder."}
{"text": "1. Field of the Invention\nThe present invention relates to wireless communication, and more particularly, to a method for transmitting a data unit in a wireless local area network system and a device for supporting the same.\n2. Related Art\nIn recent years, with the development of information and communication technology, various wireless communication technologies have been developed. Among them, a Wireless Local Area Network (WLAN) is a technology that enables a portable terminal such as a Personal Digital Assistant (PDA), a laptop computer, and a Portable Multimedia Player (PMP) to access an Internet in a wireless scheme at a house, a business, or a specific service providing zone.\nUnlike an existing wireless LAND system for supporting High Throughput (HT) and High Throughput (VHT) using 20/40/80/160/80+80 MHz bandwidth of 2 GHz and/or 5 GHz band, a wireless LAN system capable of being operated at a band less than 1 GHz is suggested. If the wireless LAN system is operated at a band less than 1 GHz, service coverage by an access point AP may be expanded as compared with an existing LAN system. Accordingly, one AP manages more STAs.\nMeanwhile, according to variation in a frequency band and a bandwidth of a used wireless channel, and rapid increase of service coverage due to this, various implementation examples with respect to a format of a new data unit usable in a next generation wireless LAN system and a transmitting method according thereto have been provided. According to the varied wireless environment and introduction of a varied format of the data unit, there is a need to suggest a method of transmitting data units capable of reducing performance degradation of a wireless LAN system and providing more efficient data processing performance."}
{"text": "Light transmissions used for communication are extremely secure due to the fact that the light transmission is focused within a narrow beam, requiring placement of equipment within the beam itself to establish a communication link. Also, because light transmissions in the visible spectrum are not regulated by the FCC, light transmissions may be used for communications purposes without the need of a license. Light transmissions are also not susceptible to interference nor do they produce noise that may interfere with other devices.\nLight emitting diodes (LEDs) may be used as light sources for data transmission, as described in U.S. Pat. Nos. 6,879,263 and 7,046,160, the entire contents of each being expressly incorporated herein by reference. LED technology provides a practical opportunity to combine lighting and communication. This combination of lighting and communication allows ubiquitous light sources to be converted to, or supplemented with, LED technology to provide for communications while simultaneously producing light for illumination purposes.\nRegarding office buildings, building management is a complex science which incorporates and governs all facets of human, mechanical and structural systems associated with buildings. As a result of the complexity, most commercial buildings are managed by commercial property management companies with great expertise. Both at the time of construction and throughout the life-cycle of a building, the interrelationships between people and the mechanical and structural systems are most desirably evaluated.\nAnother very important consideration associated with building management is energy management. Energy management is quite challenging to design into a building, because many human variables come into play within different areas within a building structure. Different occupants will have different preferences and habits. One occupant may require full illumination for that occupant to operate efficiently or safely within a space, while a second occupant might only require a small amount or local area of illumination. Further complicating the matter of energy management is the fact that many commercial establishments may experience rates based upon peak usage. A business with a large number of lights that are controlled with a common switch may have peak demands which are large as compared to total consumption of power, simply due to the amount of power that will rush into the circuit. Breaking the circuit into several switches may not adequately address inrush current, since a user may switch more than one switch at a time, such as by sliding a hand across several switches at once. Additionally, during momentary or short-term power outages, the start-up of electrical devices by the power company is known to cause many problems, sometimes harming either customer equipment or power company devices.\nEnergy management may also include consideration for differences in temperature preferred by different occupants or for different activities. Heating, Ventilation, and Air Conditioning (HVAC) demand or need is dependent not only upon the desired temperature for a particular occupant, but also upon the number of occupants within a relatively limited space\nWith careful facility design, considerable electrical and thermal energy can be saved. Proper management of electrical resources affects every industry, including both tenants and building owners. In many instances facility design has been limited to selection of very simple or basic switches, and thermostats, and particular lights, all fixed at the time of design, construction or installation.\nModern communications systems interconnect various electrical, electro-mechanical, or electrically controlled apparatuses. These connections may be referred to as connections between client devices and host devices. Host devices are simply parts of a network that serve to host or enable communications between various client devices. Generally speaking, host devices are apparatuses that are dedicated to providing or enabling communications. Peer-to-peer networks may also exist wherein, at any given moment, a device may be either client or host. In such a network, when the device is providing communication, data, information or services, it may be referred to as the host, and when the same device is requesting information, it may be referred to as the client.\nClient devices may commonly include computing devices of all sorts, ranging from hand-held devices such as Personal Digital Assistants (PDAs) to massive mainframe computers, and including Personal Computers (PCs). However, over time many more devices have been enabled for connection to network hosts, including for exemplary purposes printers, network storage devices, cameras, other security and safety devices, appliances, HVAC systems, manufacturing machinery, smart phones, mobile applications and so forth. Essentially, any device which incorporates or can be made to incorporate sufficient electronic circuitry may be so linked as a client to a host.\nMost current communications systems rely upon wires and/or radio waves to link clients and hosts. Existing client devices are frequently designed to connect to host network access points through wired connections, fiber optic connections, or as wireless connections, such as wireless routers or wireless access points.\nBuildings frequently incorporate wireless networks which are subject to a number of limitations. One of these is the lack of specific localization of a signal and device. For exemplary purposes, even a weak Radio-Frequency (RF) transceiver, in order to communicate reliably with all devices within a room, will have a signal pattern that will undoubtedly cross into adjacent rooms. When many rooms are to be covered by different transceivers, signal overlap between transceivers requires more complex communications systems, including incorporating techniques such as access control and device selection based upon identification. Radio frequency systems are subject to outside tapping and corruption, since containment of the signal is practically impossible for most buildings.\nIn addition to data communications, buildings and other spaces may also have a number of needs including, for exemplary illumination, fire and smoke detection, temperature control, and public address to name a few. With regard to illumination, buildings and other spaces are designed with a particular number and placement of particular types of light bulbs. Most designers incorporate incandescent or fluorescent bulbs to provide a desirable illumination within a space. The number and placement of these bulbs is most commonly based upon the intended use of the space.\nThe art referred to and/or described above is not intended to constitute an admission that any patent, publication or other information referred to herein is “prior art” with respect to this invention. In addition, this section should not be construed to mean that a search has been made or that no other pertinent information as defined in 37 C.F.R. § 1.56(a) exists.\nAll U.S. patents and applications and all other published documents mentioned anywhere in this application are incorporated herein by reference in their entirety.\nWithout limiting the scope of the invention, a brief summary of some of the claimed embodiments of the invention is set forth below. Additional details of the summarized embodiments of the invention and/or additional embodiments of the invention may be found in the Detailed Description of the Invention below. A brief abstract of the technical disclosure in the specification is provided for the purposes of complying with 37 C.F.R. § 1.72."}
{"text": "(a) Field of the Invention\nThe present invention relates to a waveguide type semiconductor photodetector and, more particular, to a wide-band semiconductor photodetector having a waveguide and exhibiting a high optical sensitivity and a low output distortion.\n(b) Description of the Related Art\nA conventional semiconductor photodetector incorporating therein a waveguide for guiding incident light to an optical absorption layer (hereinafter called waveguide type semiconductor photodetector or simply photodetector) has a layer structure wherein p- and n-type semiconductor layers sandwich therebetween a lightly-doped optical absorption layer or core layer to form a p-n junction. This type of semiconductor photodetector is generally applied with a reverse bias voltage between the p- and n-type semiconductor layers to deplete the optical absorption layer of carriers, and takes advantage of the high electric field generated in the depleted layer in the optical absorption layer to effect a photo-electric conversion of a signal light received through the incidence facet of the waveguide. In this process, the semiconductor photodetector receives the signal light through the incidence facet to guide the same to the optical absorption layer, and detects excited carriers generated by the incident light in the depletion layer as a photo-current. The excited carriers generated and drifting in the depletion layer are separated by the high electric field in the depletion layer as separate holes and electrons. The separated electrons reach p-type cladding layer and the separated holes reach n-type cladding layer, thereby contributing the generation of the photo-current.\nThe waveguide type semiconductor photodetector has several advantages including selection ability by the wavelength of the signal light, a high operational speed, and a wide-band characteristic. It also has the advantage of having a profile and a structure similar to those of a semiconductor laser etc. to facilitate integration therewith.\nHowever, it is difficult to adapt the mode field diameter of the conventional photodetector to the core diameter (or mode field diameter) of an optical fiber to be coupled with the photodetector for optical transmission, thereby raising a problem of coupling loss therebetween.\nThe optical adaptation in the mode field diameter between two systems, such as between the optical fiber and photodetector is discussed herein in view of the importance in reduction of the optical loss. Assuming that axial deviation between the two systems is in one direction, i.e., either horizontal or vertical direction, the coupling factor .eta. between the two systems can be represented by the following equation: EQU .eta.=2.multidot.W.sub.1 .multidot.W.sub.2 /{(W.sub.1.sup.2 +W.sub.2.sup.2).multidot.exp [-2.delta..sup.2 /(W.sub.1.sup.2 +W.sub.2.sup.2)]} (1)\nwherein W1 and W2 represent mode field diameters of both the systems, and .delta. represents the deviation or offset amount between the optical axes of both the system. The equation (1) is obtained from overlapping integral of electric field in the one direction.\nThe coupling loss between the two systems can be calculated by the above equation (1) from a specified mode field diameter ratio (or relative mode field diameter) between the two systems: if the ratio is 2, then the calculated coupling loss is 1 dB; and if the ratio is 3, then the calculated coupling loss is 2.5 dB.\nIn the case where the internal quantum effect is assumed 100% and if the mode filed diameter ratio between the two systems is 2, the optical sensitivity of the photodetector for a signal light having a 1.3 .mu.m wavelength is approximately 0.85 A/W, and if the ratio is 3 in a similar condition, the optical sensitivity is approximately 0.6 A/W.\nIt is to be noted the mode field diameter of the photodetector should be equal to that of the optical fiber or quartz waveguide to be coupled therewith in order to obtain a satisfactory coupling efficiency or to reduce the coupling loss. For this purpose, the optical absorption layer of the photodetector should be grown to have a sufficient thickness, for example, 4 .mu.m or above to be adapted with the diameter of the light emitted from the optical fiber or quartz waveguide.\nHowever, it is difficult to obtain a thickness of 4 .mu.m or above for an optical absorption layer in a photodetector wherein the layer structures including the optical absorption layer is epitaxially grown on a substrate. Namely, the adaptation of the mode field diameter by increasing the thickness of the epitaxial optical absorption layer is not practical in the photodetector.\nJP-A-4(1992)-241272 proposes a photodetector wherein the optical absorption be 0.15 .mu.m thick to reduce the confinement efficiency thereof, thereby increasing the effective mode field diameter of the optical absorption layer. This proposal, however, has a problem in which a high electric field generated within the optical absorption layer causes a zener break-down due to a tunnel current."}
{"text": "The invention relates to a method and system for determining movement of a motor in a vehicle power window system for detecting an object caught between the window and its respective frame.\nMany vehicles today have electronically controlled windows and electronically controlled sun/moon roofs. These systems provide the operator with ease in opening and closing the windows. However, if the operator is distracted while closing the window, it is possible for an object, such as an arm, hand or finger, to be caught between the window and the window frame of the automotive window.\nPower window systems typically include a regulator attached to the window for opening and closing the window and driven by an electric drive motor. The motor current wave form of these windows has a distinct sinusoidal component caused by the commutation process. Each commutation pulse can be related to angular rotation of the motor armature. Thus, window position can be determined by counting the commutation pulses. However, a problem arises due to random error pulses mixed in with the commutation pulses thereby introducing error into the position count. The error pulses cannot be filtered out using conventional band pass filter techniques because the magnitude and frequency of these error pulses are within the frequency range of the commutation pulses.\nThus, there exists a need for accurately determining movement of the motor via a change in the position of the motor for determining the presence of an object caught between a power window and its respective frame.\nA method and system is disclosed for determining movement of a motor in a vehicle power window system. A current sensor senses a current of the motor and generates an input signal having a phase associated therewith. A signal generator generates a voltage controlled reference signal having a predetermined phase and a frequency associated therewith. A comparator is in communication with the current sensor and the signal generator for comparing the phase of the input signal with the phase of the reference signal. A control circuit, in communication with the signal generator, determines whether there is any movement of the motor based on the comparison.\nThe comparator determines a difference between the phase of the input signal and the phase of the reference signal and passes this difference to the signal generator, which then adjusts the frequency of the reference signal based on this difference. A second comparator converts the adjusted reference signal into a digital signal having pulses so that the control circuit can count the number of pulses in the adjusted reference signal to determine the position of the motor. The control circuit also determines whether the position is changing. If the position does not change, this may be indicative of an obstruction.\nA motor drive circuit that controls the power window system is controlled to either stop or reverse the direction of the motor if the position of the motor has not changed and is constant.\nThese and other features of the present invention can be understood from the following specification and drawings."}
{"text": "1. Field of the Invention\nThe present disclosure relates generally to volatile material dispensing devices and operating methodologies therefore and, more particularly, to such devices and methodologies that dispense multiple volatile materials.\n2. Description of the Background\nVolatile material dispensing devices come in a variety of different forms. Some dispensing devices require only ambient airflow to disperse a liquid volatile material therefrom, e.g., from a wick extending from a volatile material container. Other devices are battery-powered or receive household power via a plug extending from the device. Some such battery-powered devices include a heating element for heating a volatile material to promote vaporization thereof. Other devices employ a fan or blower to generate airflow to direct volatile material out of the device into the surrounding environment. Still other devices that dispense volatile materials utilize ultrasonic means to dispense the volatile materials therefrom. In yet another example, some dispensing devices are configured to automatically actuate an aerosol container containing a pressurized fluid to dispense the fluid therefrom.\nIn the past, various means have been utilized to dispense one or more volatile materials from a single device. Multiple volatile materials have been used, for example, to prevent habituation, which is a phenomenon that occurs when a person becomes used to a particular volatile material such that they no longer perceive that volatile material. Alternatively or in conjunction, multiple volatile materials have been used to provide environmental effects that can be customized by a user, e.g., to provide a first fragrance in the morning to gently encourage a user to awake from sleep and a second fragrance in the evening to calm the user before falling asleep.\nDue, in part, to the variety of user preferences and needs for creating individualized environmental effects, there is an ever growing need for different volatile material dispensing devices to suit different users. Consequently, the present disclosure provides volatile material dispensers with different operating methodologies that may be preferred by some users over other devices."}
{"text": "The present invention relates to an apparatus for grinding and polishing a seam and more specifically to apparatus for grinding and polishing a welded seam between plates of stainless steel. Prior apparatus for such grinding and polishing generally consisted of a manually manipulated grinder and a separate manually manipulated polishing mechanism. Such apparatus required at least semi-skilled labor for its operation and compared to the apparatus of the present invention required a substantially greater length of time to complete the grinding and polishing operation. The apparatus of the present invention accomplishes the grinding and polishing function rapidly and effectively with a minimum of man power and with the use of operators having lesser skills than were required with prior devices."}
{"text": "When a borehole is drilled for hydrocarbon exploration, it is conventional to insert a casing string into the borehole to protect the borehole formation. A liner string can then be suspended within the casing string and can be connected to the top side by a drill string. Normally, cement is injected into the annulus between the outer surface of the casing string and the inner surface of the borehole in order to secure the casing string.\nThere are devices available that permit pressure testing of the casing string, or if present also permit pressure testing of the liner string, or which permit activation of pressure activated tools in the casing or liner string, the majority of these devices permitting pressure testing, or pressure activation respectively, after the cement has been inserted into the annulus.\nHowever, in order to be able to pressure test the casing string after cementing, it is known to run in a packer tool. The packer tool comprises an outer expandable seal that when expanded seals against the inside of the casing string, which then permits pressure testing above the site of the seal. However, using such a packer tool can be detrimental to the casing string, as the packer tool exerts very high loads on casing when pressure testing, typically in the region of 10,000-15,000 psi. Further, the packer tool must be retrieved from the well after the testing operation has been completed.\nAlternatively, a seat is provided at a suitable point on the inside of the casing or liner string so that when a plug is released, it travels down the casing or liner string and will hopefully land on the seat, thereby forming a seal so that pressure testing can occur above the plug. However, the plug and seat arrangement has the disadvantage that it is not certain that the plug will correctly land on the seat. The plug is normally released during the cementing operation and often does not land correctly on the seat, making it impossible to perform the pressure test. Further, the plug and seat pressure testing arrangement has the disadvantage that the cement is usually in position and has set or hardened by the time the pressure test is conducted. Therefore, if there is a leak in the system, the casing cannot be retrieved, resulting in expensive and time consuming remedial work to ensure pressure integrity.\nFurther, it has been known for the plug and seat pressure testing arrangement to fail during a pressure test. If this occurs, then the build up of high fluid pressure that precedes the failure can expel the cement from its intended location, and thus causes a poor cement job that requires remedial work.\nIf a packer, or the plug and seat arrangement is used after the cement has set, the high fluid pressure that is exerted on the casing can cause the metal casing to be bowed outward or expanded. This causes the cement to be displaced. Therefore, after the fluid pressure has been removed, the profile of relatively elastic metal casing will return to its original pre-pressure test state, but the cement may have set or hardened, and thus will not return to its original pre-pressure test state. Therefore, a micro-annulus may be formed between the outer diameter of the casing and the cement bond in the borehole, which may lead to gas migration up the borehole, and/or loss of zonal isolation.\nIn order to activate pressure operated tools located downhole, such as a conventional liner hanger system, or a conventional running tool for running a liner hanger system downhole, it is known to drop a ball down the casing string in the fluid path. The ball eventually lands on a ball seat located below the tool to be activated. Thus, the fluid path is blocked. This results in an increase in the fluid pressure, which can then activate the pressure operated tool. Then, when the pressure operated tool has been activated, the fluid pressure is increased such that the ball and/or the ball seat is sheared out of position, down the string. Fluid circulation can then continue in the same manner as before the ball was dropped.\nHowever, there are various problems associated with this conventional apparatus and method for activating a pressure operated tool. The time taken for the ball to reach the ball seat can be considerable, and there can be problems with getting the ball to land on the ball seat, particularly in highly deviated wells such as horizontal wells, where the ball may have to travel a relatively long distance through a horizontal section of the well. Also, when the ball and/or the ball seat has been sheared out of position, the formation receives a hydraulic shock, which can lead to a loss of circulation of the fluid."}
{"text": "The present invention is directed to a method and apparatus for changing the imaging scale in x-ray lithography. The apparatus includes a source for generating a collimated beam of radiation, means for positioning a mask in the beam of radiation before an object to be structured, an adjustment or mounting unit for positioning the object in the beam and for alignment of the object relative to the mask.\nThe progressive miniaturization of micro-electronic components places an extremely high demand on the performance capability of the lithographic methods. Thus, it is currently possible to routinely generate structures having dimension in a micrometer range (d=2-4 .mu.m), with a light-optical projection predominantly utilized in very large scale integration (VLSI) fabrication.\nIt has been suggested, as a further improvement in light optical methods, to utilize short-wave ultraviolet light having a wavelength .tau..apprxeq.200-300 nm. However, the utilization of very short-wave ultraviolet light has a lot of technical problems so that the theoretical limit of resolution of about 0.5-0.8 .mu.m can probably not be achieved.\nOne is, therefore, forced to develop new lithographic methods for producing structures in the sub-micron region. For example, see an article by H. Schaumburg \"Neue Lithografieverfahren in der Halbleitertechnik\", Elektronik 1978, No. 11, pp. 59-66. X-ray lithographic methods have, therefore, achieved special significance and their resolution is not limited by diffraction effects as a consequence of the short wavelength of the radiation, which wavelength is approximately .tau..apprxeq.0.5-4 nm, but by the range of electrons in the photoresist emitted from the layer to be structured. X-ray lithographic equipment having a conventional radiation source for a whole-wafer exposure of wafers are disclosed, for example, in an article by J. Lyman, \"Lithography steps ahead to meet VLSI challenge\", Electronics, July 1983, pp. 121-28. In these apparatus, the transfer of the prescribed structure onto the semiconductor wafer occurs on the basis of a shadow imaging in that the adjustment mask-wafer pairs are exposed with an x-radiation coming from a nearly punctiform source. The imaging of the mask structure onto the wafer surface corresponding to a conical projection occurs with a magnification scale M=1: (1+P/L), which is defined by the distance P (P.apprxeq.30 .mu.m) between the mask and wafer and the distance L (L.apprxeq.30 cm) between the x-ray source and the wafer. The changes in the rated size of the mask and the wafer, which changes occur during the manufacturing process as a consequence of thermal expansion and warping, can be compensated in a simple way in that the conical projection by the imaging scale is correspondingly adapted by changing what is referred to as the proximity distance P. This known method, however, fails when electron synchrotons or, respectively, electron storage rings are utilized as high-intensity x-ray sources. As a consequence of the high collimation degree of the synchroton radiation emitted by the electrons circulating on the circular path, the exposure of the mask-wafer pair arranged at a distance of several meters from the source or storage ring occurs on the basis of nearly exact parallel projection and, thus, L.apprxeq..infin.. In order to also guarantee a high overlay degree in the synchroton lithography, it must be assured that size variations of mask and wafer can be compensated by adapting the imaging scale."}
{"text": "1. Field of the Invention\nThis invention relates generally to safety devices for revolvers and, more specifically, to an indicia for marking a single empty chamber in an otherwise loaded revolver.\n2. Description of the Prior Art\nGun safety has long been of primary importance to gun enthusiasts. The many deaths and injuries which occur annually due to the accidental discharge of a bullet from a revolver serve to emphasize the importance of gun safety.\nIn the past, gun safety was principally geared toward mechanical safety mechanisms which would work in conjunction with the revolvers. For example, U.S. Pat. Nos. 1,842,847, 3,085,360, and 3,208,176 essentially disclosed a mechanism or means for physically blocking or preventing the undesired discharge of a bullet from a fire arm. U.S. Pat. No. 2,100,273 disclosed a mechanical means for indicating when a gun barrel was loaded.\nThese safety mechanisms were somewhat effective in preventing the accidental discharge of a bullet from a revolver. However, when these mechanisms malfunctioned, either the safety mechanism or the revolver was rendered inoperative. Furthermore these mechanisms were often clumsy and/or inconvenient to employ, causing many gun enthusiasts to use the mechanisms less than they otherwise would have.\nThe most effective method for assuring safe handling of a gun was simply to keep the revolver unloaded. However, this method was totally inapplicable to situations in which the revolver was about to be used or situations where it was essential to fire the revolver quickly once the need or desire to do so was perceived. In the May 1981 issue of Guns and Ammo magazine, it was stated that the best method for gun safety under these circumstances was to carry the revolver with the hammer down on an empty chamber.\nU.S. Pat. No. 3,407,526 disclosed a revolver capable of firing a variety of loads or bullet types and having an indicia on the cylinder of the revolver for indicating what types of loads are in each chamber of the revolver. However, these indicia were limited to identification of chambers compatable with particular load types for use with that particular type of revolver and did not in any way fill the need for indicating a particular empty chamber for use of a standard revolver consistant with the practice of keeping the hammer down on an empty chamber.\nThus a need exists to provide means for rapidly and conveniently identifying the specific empty chamber onto which the hammer was to rest so as to prevent the underside discharge of a bullet therefrom."}
{"text": "The subject invention generally relates to a curable, powder-based coating composition for coating a substrate. More specifically, the subject invention relates to a curable, powder-based coating composition that includes a powder-based binder and a color effect-providing pigment for coating a substrate that has a first color effect with a film layer. Upon application of the film layer to the substrate, the color effect-providing pigment interacts with the first color effect of the substrate to produce a second color effect that is different from the first color effect of the substrate.\nPowder-based coating compositions are known in the art. A film layer of such powder-based coating compositions is applied to a substrate throughout various industries, such as the automotive coating industry, for certain functional and aesthetic purposes. U.S. Pat. Nos. 5,379,947; 5,552,487; 5,569,539; 5,601,878; and 5,639,821 all disclose various powder-based coating compositions known in the prior art. It is also known in the art to incorporate conventional pigments, and even conventional flake pigments, specifically mica or aluminum flake pigments, into powder-based coating compositions.\nThe powder-based coating compositions of the prior art, even those conventional coating compositions that incorporate conventional mica or aluminum flake pigments, are deficient because, upon application of the film layer of the coating composition to the substrate, the film layers of the prior art compositions do not achieve a suitable color effect that varies from an original color effect of the substrate. These prior art compositions cannot attain the suitable color effect because the various pigments incorporated into the powder-based coating compositions, even mica and aluminum flakes, do not appropriately interact with light waves to establish angle-dependent color and lightness effects that are responsible for achieving the suitable color effects. As such, coating systems that rely on such conventional powder-based coating compositions first require application of a color-providing basecoat layer, or other coating layer, that underlies the film layer of the powder-based coating composition to provide the angle-dependent color and lightness effects.\nIn the interest of eliminating the necessity for any color-providing basecoat layer, and also because the film layers of the prior art compositions do not achieve suitable color effects that vary from the original color effect of the substrate, it is desirable to implement a unique powder-based coating composition that incorporates a color effect-providing pigment interacting with a first color effect of a substrate to produce a second color effect differing from the first color effect. It is also desirable that the powder-based coating composition of the subject invention does not require application of an underlying color-providing basecoat layer to achieve the second color effect.\nA curable, powder-based coating composition is disclosed. The coating composition is utilized for coating a substrate, having a first color effect, with a film layer of the coating composition. Application of the film layer to the substrate produces a second color effect different from the first color effect of the substrate.\nThe coating composition includes a powder-based binder. The powder-based binder is the reaction product of a resin and a cross-linking agent. More specifically, the resin has a functional group and the cross-linking agent is reactive with the functional group of the resin. The coating composition also includes a color effect-providing pigment. The color effect-providing pigment includes a pigment substrate and an inorganic coating. More specifically, the pigment substrate has first and second substantially parallel and planar surfaces, and the inorganic coating is disposed on at least one of the first and second substantially parallel and planar surfaces of the pigment substrate. Furthermore, the inorganic coating disposed on the pigment substrate has an index of refraction of 1.8 or less.\nApplication of the film layer of the powder-based coating composition to the substrate allows the inorganic coating, having the index of refraction of 1.8 or less, and the pigment substrate of the color effect-providing pigment to interact with the first color effect of the substrate to produce the second color effect. The subject invention, therefore, provides a unique powder-based coating composition for coating a substrate that incorporates particular color effect-providing pigments to interact with the first color effect of the substrate to produce the second color effect. Accordingly, the coating composition of the subject invention also allows for the elimination of any color-providing basecoat underlying the film layer of the powder-based coating composition.\nThe curable, powder-based coating composition of the subject invention coats at least one surface of a substrate, having a first color effect, with a film layer. It is to be understood that the powder-based coating composition of the subject invention includes exclusively powder-based coating compositions as well as powder slurry-based coating compositions. The first color effect is the original color and original appearance of the substrate. Application of the film layer of the powder-based coating composition to the substrate produces a second color effect that is different from the first color effect of the substrate.\nAlthough metallic substrates, such as automotive body panels, are typical, the powder-based coating composition may be applied to other substrates without varying the scope of the subject invention. By way of example, the powder-based coating composition may be applied to a plastic substrate. Also, the powder-based coating composition of the subject invention is primarily utilized as a powder clearcoat applied to the substrate to produce the second color effect without an underlying color-providing basecoat film layer. The powder-based coating composition of the subject invention may also be utilized with an underlying film layer, such as the underlying color-providing basecoat film layer, where the underlying film layer is actually the substrate to which the film layer of the powder-based coating composition is applied. In other words, the substrate is not required to be a bare automotive body panel.\nThe powder-based coating composition includes a powder-based binder and a color effect-providing pigment. The color effect-providing pigment is described below. The powder-based binder is a film-forming binder that is the reaction product of a resin and a cross-linking agent. The resin includes a functional group, and the cross-linking agent is specifically reactive with the functional group of the resin. More specifically, the resin of the power-based binder is selected from the group consisting of acrylic resins, epoxy resins, phenolic resins, polyester resins, urethane resins, and combinations thereof. The functional group of the resin is selected from the group consisting of epoxy functional groups, carboxy functional groups, hydroxy functional groups, and combinations thereof. The cross-linking agent reactive with the functional group of the resin is selected from the group consisting of aminoplasts, blocked isocyanates, polycarboxylic acids, acid anhydrides, polyamines, and combinations thereof.\nThe color effect-providing pigment includes a pigment substrate and an inorganic coating. The subject invention preferably combines from 0.1 to 10 parts by weight of the color effect-providing pigment based on 100 parts by weight of the powder-based binder. The pigment substrate has first and second substantially parallel and planar surfaces, and the inorganic coating is disposed or applied on at least one of the first and second substantially parallel and planar surfaces of the pigment substrate. Preferably, the inorganic coating is disposed or applied on both the first and second substantially parallel and planar surfaces. The pigment substrate is preferably a platelet-shaped pigment substrate. In the context of the subject invention, the terminology platelet-shaped indicates that the pigment substrate is a minute, flattened body. Furthermore, the pigment substrate is preferably selected from the group consisting of metallic pigment substrates, non-metallic pigment substrates, and combinations thereof, depending on the particular embodiment of the subject invention. The inorganic coating also has an index of refraction of 1.8 or less. The inorganic coating and the pigment substrate, and other optional coatings as set forth below, establish a symmetrical, multilayer interference structure of the color effect-providing pigment.\nThe inorganic coating of the color effect-providing pigment varies depending on the embodiment. A suitable example for the inorganic coating is an inorganic coating including a metal oxide. The inorganic coating may also be selected from the group consisting of metal oxides, magnesium fluoride, and combinations thereof. Further suitable examples for the inorganic coating of the color effect-providing pigment include inorganic coatings selected from the group consisting of silicon oxide, silicon oxide hydrate, aluminum oxide, aluminum oxide hydrate, titanium oxide, titanium oxide hydrate, zinc sulfide, magnesium fluoride, and combinations thereof.\nAs described below, the inorganic coating and the pigment substrate of the color effect-providing pigment interact with the first color effect of the substrate to produce the second color effect upon application of the film layer of the powder-based coating composition to the substrate. When the color effect-providing pigment is incorporated into the powder-based binder according to the subject invention, the interaction of the inorganic coating and the pigment substrate with the first color effect to produce the second color effect is further defined as interference of light waves. In this embodiment, the interference of the light waves establishes angle-dependent color and lightness effects to achieve the second color effect. The interaction of the inorganic coating and the pigment substrate may also be defined as absorption of light waves to establish the angle-dependent color and lightness effect to achieve the second color effect, or as reflection of light waves to establish the angle-dependent color and lightness effects to achieve the second color effect. In one embodiment of the subject invention, the inorganic coating and the pigment substrate interact with the first color effect of the substrate such that the second color effect is different from the first color effect at least by xcex94L 20.0, xcex94a 10.0, and xcex94b 15.0 as measured according to CIELab color space.\nThe color effect-providing pigment further includes a reflective, absorbing coating which is at least partially transparent to visible light. For descriptive purposes of the subject invention, use of xe2x80x9cat least partially transparent to visible light,xe2x80x9d throughout the description indicates that the pigment substrate, the inorganic coating, or other coatings that are described in such terms, such as the reflective, absorbing coating introduced immediately above, generally transmit at least 10%, preferably at least 30%, of incident light. The reflective, absorbing coating includes a selectively absorbing metal oxide, or a non-selectively absorbing metal, or both. For descriptive purposes of the subject invention, the terminology metal oxide, as used herein, is also intended to encompass metal dioxides, metal trioxides, and so on. The reflective, absorbing coating is preferably disposed or applied on the inorganic coating.\nThe color effect-providing pigment further comprises an outer coating disposed or applied on the reflective, absorbing coating. The outer coating is different from the reflective, absorbing coating and preferably includes a selectively absorbing metal oxide. The symmetrical, multilayer interference structure includes the pigment substrate, the inorganic coating, the reflective, absorbing coating, and the outer coating.\nIn an embodiment where the pigment substrate of the color effect-providing pigment is a metallic pigment substrate, the most preferred metallic pigment substrate is aluminum. Other suitable metallic pigment substrates include, but are not limited to, all metals and alloys in platelet form known as metallic pigment substrates, such as steel, copper, copper alloys including brass and bronze, and aluminum bronze. The aluminum pigment substrate may be a passivated or an unpassivated aluminum pigment substrate. The aluminum pigment substrate preferably has an average particle size of from 5 to 50, preferably from 10 to 20, and most preferably from 13 to 16, microns. Alternatively, the aluminum pigment substrate may have a particle-size distribution where 50% of the aluminum pigment substrate has a particle size of from 13 to 16 microns and where no more than 5% of the aluminum pigment substrate has a particle size of greater than 50 microns.\nThe inorganic coating disposed on the aluminum pigment substrate preferably comprises a metal oxide. Alternatively, the inorganic coating disposed on the aluminum pigment substrate may be selected from the group consisting of silicon oxide, silicon oxide hydrate, aluminum oxide, aluminum oxide hydrate, and combinations thereof. In either embodiment for the inorganic coating, the inorganic coating has an index of refraction of 1.8 or less, preferably 1.6 or less. Also in either embodiment for the inorganic coating, the inorganic coating disposed on the aluminum pigment has a thickness of from 200 to 600, preferably from 300 to 500, nanometers (nm). It is to be understood that the thickness of the inorganic coating, and all other coatings described in the subject invention, varies as a function of the properties of the components selected for the inorganic coating. For instance, the thickness of an inorganic coating including silicon oxide may differ from the thickness of an inorganic coating including aluminum oxide.\nWhere the pigment substrate is the aluminum pigment substrate, the color effect-providing pigment optionally further comprises a reflective, selectively absorbing metal oxide. The metal oxide is disposed on the inorganic coating and has an index of refraction of 2.0 or greater and is at least partially transparent to visible light. In such embodiments, the index of refraction of the reflective, selectively absorbing metal oxide is more preferably 2.4 or greater. If present, the reflective, selectively absorbing metal oxide preferably has a thickness of from 1 to 500, more preferably from 10 to 150, nm.\nThe color effect-providing pigment in this embodiment optionally further comprises an absorbing, outer coating. The absorbing, outer coating is different from the reflective, selectively absorbing metal oxide. Furthermore, the absorbing, outer coating is disposed on the reflective, selectively absorbing metal oxide. Preferably, the absorbing, outer coating is selected from the group of selectively absorbing oxides consisting of iron (III) oxide, chromium (III) oxide, vanadium (V) oxide, titanium (III) oxide, and combinations thereof. Alternatively, the absorbing, outer coating may be selected from the group of non-selectively absorbing oxides consisting of titanium dioxide, zirconium oxide, and combinations thereof. If present, the absorbing, outer coating has a thickness of from 1 to 200, more preferably from 10 to 150, nm.\nThe symmetrical, multilayer interference structure of the color effect-providing pigment, where the metallic pigment substrate is the aluminum pigment substrate, includes Fe2O3 as the reflective, selectively absorbing metal oxide, SiO2 as the inorganic coating, Al as the metallic pigment substrate, SiO2 as the inorganic coating, and Fe2O3 as the reflective, selectively absorbing metal oxide. That is, this color effect-providing pigment has a symmetrical, multilayer interference structure of Fe2O3/SiO2/Al/SiO2/Fe2O3. Such color effect-providing pigments having the metallic pigment substrate are commercially available from BASF Corporation, Southfield, Mich. as Variocrom(copyright) Magic Red K 4411 (formerly ED 1479) and Magic Gold K 1411, and are set forth in U.S. Pat. No. 5,607,504, the disclosure of which is incorporated herein by reference in its entirety. With Variocrom(copyright) Magic Red K 4411, the second color effect is produced as a result of a color shift from red-to-yellow. With Variocrom(copyright) Magic Gold K 1411, the second color effect is produced as a result of a color shift from greenish gold-to-reddish gray. It is to be understood that the color shifts that produce the second color effect are primarily driven by the thickness of the SiO2 inorganic coating.\nThe metallic pigment substrate may also be selected from the group consisting of chromium, nickel, and combinations thereof. If the metallic pigment substrate is chromium or nickel, the color effect-providing pigment, as a whole, has an average particle size of from 5 to 40, preferably from 20 to 40, microns. Alternatively, when the metallic pigment substrate is chromium or nickel, the color effect-providing pigment, as a whole, can have a particle size distribution where no more than 10% of the pigment has a particle size of greater than 50 microns and substantially none of the pigment has a particle size of greater than 125 microns.\nIn this embodiment, the inorganic coating disposed on the metallic pigment substrate is a dielectric inorganic coating having an index of refraction of 1.65 or less. The inorganic coating having the index of refraction of 1.65 or less is selected from the group consisting of silicon oxide, silicon oxide hydrate, aluminum oxide, aluminum oxide hydrate, magnesium fluoride, and combinations thereof.\nIn this embodiment, the color effect-providing pigment optionally further includes a semi-transparent metal coating disposed on the inorganic coating. The semi-transparent metal coating most preferably includes aluminum. Alternatively, the semi-transparent metal coating is selected from the group consisting of aluminum, gold, copper, silver, and combination thereof.\nThe symmetrical, multilayer interference structure of the color effect-providing pigment, where the metallic pigment substrate is the chromium or nickel pigment substrate includes Al as the semi-transparent metal coating, SiO2 or MgF2 as the inorganic coating, and chromium or nickel as the metallic pigment substrate. Such color effect-providing pigments are commercially available from Flex Products, Inc., Santa Rosa, Calif., and are set forth in U.S. Pat. No. 5,135,812 and U.S. patent application Ser. No. 08/172,450, the disclosures of which are incorporated herein by reference in their entirety.\nThe metallic pigment substrate may alternatively be steel. In the context of the subject invention, it is to be understood that steel is an alloy of iron and from 0.02 to 1.5 parts carbon. If the metallic pigment substrate is steel, then it is most preferably stainless steel. One suitable example for more generally defining the steel pigment substrate is as an alloy of steel having from 1 to 30 parts by weight of chromium based on 100 parts by weight of the alloy of steel. In the most preferred embodiment, the metallic pigment substrate is selected from the group consisting of aluminum, chromium, nickel, steel, stainless steel, and combinations thereof.\nAs described above, the pigment substrate may be a non-metallic pigment substrate. The non-metallic pigment substrate has an index of refraction of 2.0, preferably 2.4 or greater. The non-metallic pigment substrate may be iron oxide, mica having an oxide coating, or combinations thereof. In an embodiment where the non-metallic pigment substrate is the mica having the oxide coating, the oxide coating is more specifically defined as a TiO2 coating having a thickness of from 10 to 300 nm. As with the metallic pigment substrate, the non-metallic pigment substrate has an average particle size of from 5 to 50 microns. Preferably, the average particle size of the non-metallic pigment substrate is from 10 to 30, and most preferably from 15 to 20 microns.\nWith the non-metallic pigment substrate, the inorganic coating disposed on the substrate, is preferably selected from the group consisting of metal oxides, magnesium fluoride, and combinations thereof. Alternatively, the inorganic coating disposed on the non-metallic pigment substrate is selected from the group consisting of silicon oxide, silicon oxide hydrate, aluminum oxide, aluminum oxide hydrate, and combinations thereof. The inorganic coating has a thickness of from 20 to 800, preferably from 50 to 600, nm.\nThe color effect-providing pigment optionally further includes a reflective, absorbing coating disposed on the inorganic coating. The reflective, absorbing coating is selected from the group consisting of metals, metal oxides, metal sulfides, metal nitrides, and combinations thereof. The reflective, absorbing coating has a thickness of from 1 to 500, preferably from 10 to 150, nm.\nThe color effect-providing pigment optionally further includes an absorbing, outer coating. The absorbing, outer coating is different from and is disposed on the reflective, absorbing coating described above. Preferably, the absorbing, outer coating comprises a metal oxide. Alternatively, the absorbing, outer coating may be selected from the group consisting of silicon oxide, silicon oxide hydrate, aluminum oxide, aluminum oxide hydrate, tin oxide, titanium dioxide, zirconium oxide, iron (III) oxide, chromium (III) oxide, and combinations thereof.\nThe symmetrical, multilayer interference structure of the color effect-providing pigment, where the non-metallic pigment substrate is the iron oxide pigment substrate, includes Fe2O3 as the reflective, absorbing coating, SiO2 as the inorganic coating, Fe2O3as the non-metallic pigment substrate, SiO2 as the inorganic coating, and Fe2O3 as the reflective, absorbing coating. That is, this color effect-providing pigment has a symmetrical, multilayer interference structure of Fe2O3/SiO2/Fe2O3/SiO2/Fe2O3. Such color effect-providing pigments having the non-metallic pigment substrate are commercially available from BASF Corporation, Southfield, Mich. as Variocrom(copyright) Magic Purple K 5511 (formerly ED 1480), and is set forth in U.S. Pat. No. 5,958,125, the disclosure of which is incorporated herein by reference in its entirety. With Variocrom(copyright) Magic Purple K 5511, the second color effect is produced as a result of a color shift from violet-to-gold. As above, it is to be understood that the color shift that produces the second color effect with the Variocrom(copyright) Magic Purple K 5511 is primarily driven by the thickness of the SiO2 inorganic coating.\nA coating system is also disclosed. The coating system includes the substrate, preferably the automotive body panel, having the first color effect. The coating system also includes the film layer of the powder-based coating composition as described above. The film layer is at least partially-transparent to visible light. As such, the most preferred coating system is where the powder-based coating composition is a powder clearcoat applied on the substrate to produce the second color effect.\nThe coating system may optionally include a second film layer. Preferably, the second film layer is also at least partially-transparent to visible light. The second film layer is applied on the film layer of the powder-based coating composition. The purpose of the application of the second film layer on the first film layer is primarily to enhance appearance characteristics, such as gloss, of the film layer, if necessary. For instance, the coating system may be a high-gloss coating system when the second film layer is included. More specifically, with the second film layer, the coating system utilizing the powder-based coating composition of the subject invention has a 20 degree gloss of at least 65, preferably of at least 75, as defined by ASTM D523-89 (Re-Approved 1999). Alternatively, the coating system utilizing the powder-based coating composition may have a 60 degree gloss of at least 75, preferably of at least 85, as defined by the same ASTM standard. The 20 and 60 degree glosses are preferably measured with a BYK-Gardner Micro-Gloss Meter, specifically Model No. GB-4501. Alternatively, a BYK-Gardner Haze-Gloss Meter, preferably Model Nos. GB-4601 and GB-4606, may be utilized having a different scale of gloss units. In general, if the powder-based coating composition according to the subject invention is not utilized as strictly a single-layer, powder clearcoat, then it may be utilized as a color-providing basecoat film layer with a second, non-pigmented clearcoat applied over the color-providing basecoat film layer.\nThe coating system may alternatively include an underlying film layer applied to the substrate prior to application of the film layer of the powder-based coating composition. Where the underlying film layer has been applied, the underlying film layer is the substrate to which the film layer of the powder-based coating composition is applied. As such, it is to be understood that the underlying film layer can be an electrocoat film layer, a primer surfacer film layer, or a color-providing base coat film layer as known in the art.\nA method for coating the substrate is also disclosed. Generally, the method for coating the substrate to produce the second color effect upon application of the film layer is characterized by using the powder-based coating composition set forth above. More specifically, in the subject method, the powder-based binder and the color effect-providing pigment are combined to establish the powder-based coating composition. Preferably, the powder-based binder and the color effect-providing pigment are combined in amounts from 0.1 to 10, more preferably from 1 to 6, parts by weight of the color effect-providing pigment based on 100 parts by weight of the powder-based binder.\nThe step of combining the powder-based binder and the color effect-providing pigment varies depending on the embodiment of the subject invention. In one embodiment, the step of combining is further defined as dry blending the color effect-providing pigment into the powder-based binder. In the dry blending embodiment, it may be particularly important that the pigment substrate have an increased particle size toward the upper limit of from 5 to 50, preferably from 25 to 40, xcexcm. The increased particle size of the pigment substrate in the dry blending embodiment is important to minimize pigment settling and flocculation concerns and also to ensure that the color-effect providing pigment produces the second color effect. The dry blending embodiment, further includes the step of agitating the dry blend of the color effect-providing pigment and the powder-based binder. As such, the color effect-providing pigment is uniformly dispersed throughout the powder-based binder.\nIn another embodiment, the step of combining is further defined as extruding the color effect-providing pigment into the powder-based binder. In the extruding embodiment, it may be particularly important that the pigment substrate is stainless steel such that the pigment substrate of the color effect-providing pigment can effectively withstand the forces typically involved in the extruding of the color effect-providing pigment. The extruding embodiment further comprises the step of milling the extrusion of the color effect-providing pigment and the powder-based binder to establish the powder-based coating composition.\nFinally, the step of combining according to the various preferred embodiments may also be defined as bonding, more specifically impact bonding, the color effect-providing pigment with the powder-based binder.\nThe subject method further includes the step of applying the powder-based coating composition to the substrate. Upon application of the powder-based coating composition, the second color effect is produced as a result of the interaction of the inorganic coating and the pigment substrate with the first color effect of the substrate as described above. It is to be understood that the most preferred manner in which to apply the powder-based coating composition is by spray application. Finally, the film layer of the powder-based coating composition is cured such that the film layer produces the second color effect.\nThe following examples illustrating the formation of the powder-based coating composition according to the subject invention and illustrating certain properties of the film layer of the coating composition applied on the substrate, as presented herein, are intended to illustrate and not limit the invention."}
{"text": "The computer industry is currently striving to make computers and computer application programs more \"user-friendly.\" Many user-friendly computer application programs permit users to input commands via a graphical user interface, e.g., a screen displaying various options in graphical format. Software programs using the WINDOWS.TM. environment, manufactured by Microsoft Corporation, often take advantage of the user-friendly interface permitted under this environment, by allowing users to enter commands by moving a cursor or pointer on the computer screen to select various options or move through windows.\nIn today's computers, cursor movement is often controlled using input devices such as mice or trackballs. These input devices also allow selection of options and navigation through windows and menus, however, are unsuitable for applications requiring high resolution computer input such as drawing, inking, gesturing, recognition and absolute positioning. Another drawback of mice and trackballs are that they are attached to the computer by a cord. At times, this cord can be annoying or a hindrance.\nPen and tablet and other similar digitizing computer input devices permit high resolution input. However, these input devices are not readily suited for cursor control, selection of options and navigation through windows and applications. Another drawback of pen and tablet input devices are that they require a special digitizing tablet upon which to move the pen. This tablet is often bulky and unsuitable in some working environments.\nU.S. Pat. No. 5,166,668 describes a wireless computer input system using a battery-powered pen-type input device having a light emitting element at its tip. Light generated by a light-emitting diode in a pen is received by an input unit. The input unit detects and calculates the pen's position and outputs the position coordinates to a computer. The pen must be moved within a small area proximate to the input unit. The input unit has two sets of optical elements each having: a lens, an optical detector, and an optical filter. The optical detector is described as a two-division photodiode which detects a ratio of light impinging on both halves of the detector.\nThe system described in the above patent relies on two-division photodiodes. Two-division photodiodes require both halves to be tuned to provide equal outputs in response to an equal amount of impinging light, i.e., the two divisions must be identical. Such photodiodes are expensive to produce and may suffer from less than optimal resolution. The incident light spot must be large enough to impinge on both halves of the photodiode, and with enough energy, to produce ample current at the output terminals. Therefore, the photodiode output is dependent on the amount of incident light received. This often requires the light emitting element to produce a large amount of light, quickly draining the pen's battery during frequent use.\nU.S. Pat. No. 5,045,843 describes another computer input system using a moveable, wireless remote unit and a stationary input unit, both transmitting signals between each other. The remote unit contains a light emitting element, a light detector, a switch, a power supply, and a controller. The stationary input unit contains a light emitting element, a photo detector, an amplifier, and a processor.\nOne problem with the system described in this patent is that the light detector and associated circuitry are included in a large and cumbersome wireless remote unit. The light detector is placed at the end of a long tube, forming a camera-type receiver. The receiver, and the system as a whole, would be expensive to manufacture because of the many and costly components necessary to operate the invention. Another problem is that the remote unit must be moved in a large area, and would be unsuitable for small, detailed and accurate movements on a small work surface. The remote unit is used primarily to move a cursor on a television screen at a substantial distance (several feet) from the television. Consequently, this system cannot be used for inking, gesturing, or character recognition, because these functions require precision recognition of a user's small hand movements."}
{"text": "Brachycome multifida cultivar Metallic Blue.\nThe present invention relates to a new and distinct cultivar of Brachycome plant, botanically known as Brachycome multifida and hereinafter referred to by the name xe2x80x98Metallic Bluexe2x80x99.\nThe new Brachycome was discovered by the Inventor in Croydon, Victoria, Australia as a seedling from a random cross of two unidentified selections of Brachycome multifida, not patented. The new Brachycome was discovered and selected by the Inventor as a plant within the progeny of the stated cross in a controlled environment in Croydon, Victoria, Australia in 1991. The selection of the new Brachycome was based on its upright growth habit and inflorescences with light violet-colored ray florets. Compared to plants of the parent selections, plants of the new Brachycome have a more upright plant growth habit, darker green-colored leaves and larger inflorescences.\nAsexual reproduction of the new Brachycome by cuttings taken in a controlled environment in Montrose, Victoria, Australia, since 1991, has shown that the unique features of this new Brachycome are stable and reproduced true to type in successive generations.\nThe new Brachycome has not been observed under all possible environmental conditions. The phenotype may vary somewhat with variations in environment such as temperature, daylength and light intensity without, however, any variance in genotype.\nThe following characteristics have been repeatedly observed and are determined to be basic characteristics of xe2x80x98Metallic Bluexe2x80x99 and distinguish the new Brachycome as a new and distinct cultivar:\n1. Upright and rounded plant growth habit.\n2. Large inflorescences with light violet-colored ray florets.\n3. Erect peduncles that hold inflorescences above the foliage.\nPlants of the new Brachycome can be compared to plants of the cultivar City Lights, disclosed in U.S. Plant Pat. No. 11,646. In side-by-side comparisons conducted in St. Richmond, Victoria, Australia, plants of the new Brachycome differed from plants of the cultivar City Lights in the following characteristics:\n1. Plants of the new Brachycome were more upright and taller than plants of the cultivar City Lights.\n2. Inflorescences of plants of the new Brachycome had fewer ray florets than inflorescences of plants of the cultivar City Lights.\n3. Plants of the new Brachycome and the cultivar City Lights differed in ray floret coloration.\n4. Inflorescences of plants of the new Brachycome had smaller discs than inflorescences of plants of the cultivar City Lights.\nPlants of the new Brachycome can also be compared to plants of the cultivar Billabong Mauve Delight, disclosed in U.S. Plant Pat. No. 10,889. In side-by-side comparisons conducted in St. Richmond, Victoria, Australia, plants of the new Brachycome differed from plants of the cultivar Billabong Mauve Delight in the following characteristics:\n1. Plants of the new Brachycome were more upright and taller than plants of the cultivar Billabong Mauve Delight.\n2. Inflorescences of plants of the new Brachycome were larger and had more ray florets than inflorescences of plants of the cultivar Billabong Mauve Delight.\n3. Plants of the new Brachycome and the cultivar Billabong Mauve Delight differed in ray floret coloration.\n4. Plants of the new Brachycome had longer peduncles than plants of the cultivar Billabong Mauve Delight."}
{"text": "The present invention has its application within the telecommunication sector, more specifically, relates to the analysis of mobile user traffic.\nMore particularly, the present invention refers to a method for detecting from mobile traffic applications (apps) of mobile user terminals (smartphones, tablets, etc.).\nSmartphones offer users the possibility to install on them whatever applications (apps) they decide to (apart from preinstalled apps). These apps belong to categories as entertainment, sports, productivity, travel . . . . Therefore, applications installed in a certain smartphone provide useful information about its user profile, to be understood as the set of habits and preferences of a person.\nThose apps require Internet connection for tasks as content update or access authorization. The set of queries or requests sent to retrieve data from Internet is here defined as traffic. Being a protocol a set of predefined rules that defines the way of transferring information, requests information may vary depending on the used protocol. Examples of information appearing in requests are: source IP, request date and time, domain or user agent. The latter concepts, domain and user agent, are defined as follows:                Domain is the unique name that identifies a website on the Internet.        User agent includes information about several aspects like: application source, device operating system or software version. It has to be emphasized that not every protocol includes the user agent field.        \nThe smartphones have recently experimented an exponential growth in terms of number of users and hours spent with them. In this context, knowing applications used by a customer will allow to precisely define its profile. A correct user profile is the key to success in multiple use-cases like recommender systems, protection against possible security threats (malicious apps) or statistical analysis, as defined as follows:                Recommender systems: Those systems are present in several areas such as cinema, music or shopping. They aim to predict user interests, i.e. user profile, on those areas using information of his activity. Based on these predictions, they provide recommendations to users about elements that match their interests. As more precise the predictions, the better the recommendations.        Malicious apps: Applications classified as malicious are, for example, those tricking users into unwanted pays or subscriptions.        Statistical analysis: Analysis over user profiles and their distribution, which can guide, for example, further commercial or investment decisions.        \nMobile Network Operators (MNOs) can obtain information required to define users profile from mobile traffic. A request is generated each time a mobile user interacts with an app on its smartphone. The request passes through the MNO infrastructure, which both stores it in a database as sends it to the Internet. Data stored in the MNO database is simplified information of HTTP and DNS requests. Hypertext Transfer Protocol (HTTP) is a protocol for transferring hypermedia files. Domain Name System (DNS) is a naming system for clients or services connected to the Internet or to a private network. DNS associates a domain name with an internet protocol address (IP). The information stored in the database is the domain and the date and time of the request, i.e. the complete URL is not consulted in any case. In addition, all stored data are anonymized.\nThere are approaches for analyzing mobile traffic based on domain information. However, relation between domains and applications is not bijective. Unique domains, i.e. domains exclusively accessed by an app, are the less frequent. Instead, there are some domains accessed by many apps. In the latter case, the knowledge of the domain does not univocally define the application.\nThere are also approaches for analyzing mobile traffic based on user agent. User agent presents two major drawbacks: not all HTTP petitions have user agent value, and applications developers decide the value of user agent field, so they can use another apps' user agent instead of setting their own.\nFinally, there is a great variability in the requests of a concrete application. It is due, inter alia, to the different operating systems or mobile user devices (smartphones, tablets). Even different executions of the same application on the same terminal do not maintain the same request order. Some issues related to the requests variability are, to name but a few: request may be cached, latencies between requests vary depending on the mobile use, list of domains consulted by an app may vary between devices, or dynamic content include noise in executions.\nTherefore, it is highly desirable to develop a method of apps detection from the mobile traffic which allows the MNOs to get a more precise user profile."}
{"text": "(none)\n(none)\nThe invention relates to a cervical restraint device for use on a psittacine bird to prevent feather plucking, self-inflicted trauma, and removal of or damage to medically applied healing aids used in the treatment of an injured bird. The device may be applied for short or long periods of time and may be repeatedly placed on or taken off the bird with relative ease.\nPrior art forms of avian cervical restraint collars are either constructed in a flat or in a conical structure or a combination of both. The flat or conical circular collar is worn around the neck of the bird and extending outwardly equally in all directions to act as a physical barrier to a bird attempting to reach around the collar. The cylindrical tube collar is similarly positioned around a bird\"\"s neck in partial or full extension to prevent the beak of the bird from reaching lower body parts by restricting the bending motion of the neck. Additionally, cloth padding has been added to prior art forms.\nMany birds are able to grasp these collars on the peripheral edges, at overlapping edges, or at exposed fasteners to effectively chew the collar down to an ineffective size, release the fasteners, or endanger themselves with sharp or jagged edges as a result of the chewing. Fasteners include metal snaps, nuts and bolts, plastic strips woven through slots, cloth ties, and tape.\nU.S. Pat. No. 5,197,414 illustrates a device mentioned above. This patent shows a dog collar formed as a funnel or a cone and is designated to be primarily used on dogs. The collar is used mainly to prevent a groomer from being bitten or the dog from licking wounds on its body.\nU.S. Pat. No. 5,697,328 discloses a therapeutic collar for birds that includes a sheet formed in a rectangular configuration. The sheet is then formed into a cylindrical configuration. The end edges are releasably fastened to each other. The invention at hand is formed into a sphere having a continuously curved outer surface.\nU.S. Pat. No. 5,779,828 describes a protective pet collar that is made of two sheets of flexible material and is fastened around an animal\"\"s neck for preventing the animal from bending its neck and thereby prevent the animal from biting or licking an affected injured area. The invention at hand is a rigid sphere.\nU.S. Pat. No. 5,787,842 reveals a radially projecting restrictive pet collar using an adjustable fastener to apply it to a pet\"\"s neck. The collar forms a cylindrical shape which impedes the pet from passing through narrow openings or from chewing areas of it\"\"s body under medical treatment. The invention at hand is a sphere with a curved outer surface and a smaller comparable radius.\nU.S. Pat. No. 5,915,337 refers to an adjustable tubular or cylindrical shaped cervical pet collar. The invention at hand is a nonadjustable sphere.\nU.S. Pat. No. 6,044,802 is a soft, round pillow shaped veterinary recovery collar with a central aperture and draw string to adjust to the size of the neck. The invention at hand is a nonadjustable, rigid sphere with two apertures directly opposed to one another.\nU.S. Pat. No. 6,129,054 describes a substantially rectangular padded collar which wraps around and contacts an animal\"\"s entire neck from mandible to scapula. The overlapping ends of the tubular or cylindrical collar are held in place by one or more inelastic bands. The invention at hand is an unpadded, rigid sphere with internal fasteners and or external fasteners additionally having two openings with curved borders that come in contact with only the upper and lower portions of the neck of a bird.\nThe psittacine bird has a large head relative to a long, slender neck. The neck rests in the shape of an xe2x80x9cSxe2x80x9d curve which may be straightened to increase it\"\"s overall length and reach. This allows for fitting of a semi-rigid, hollow, three dimensional cervical collar in the shape of a round, elliptical, or eccentric oval sphere, to restrict the reach of the bird\"\"s head and beak below the neck, thereby preventing self-destructive actions of feather picking or self mutilation and serving as an effective restraint to help inhibit removal of bandages or other medically applied devices. However, visual capabilities are not impaired. There is no interference with the movement of appendages.\nThe description below makes reference to a spherical form or shape. A sphere is defined as a shape which is equidistant in all directions from a central point. However, in the context of the invention at hand, the properties of the sphere as it relates to the avian cervical restraint collar, shall represent any of the three-dimensional shapes having a continuous curved surface including round, elliptical, centric and eccentric oval or multiple combination thereof.\nThe equator of a round or elliptical sphere is equidistant from either pole of the sphere. The equator of an eccentric oval sphere may be closer to one pole or the other thereby creating a top heavy or bottom heavy sphere. The sphere could even take on the general shape of an hour glass with two wide sections separated by a narrow section.\nA hollow, spherical avian cervical restraint collar is formed by engaging two halves of a durable, semi-rigid lightweight, non toxic and solid material such as but not limited to plastic. The two halves of the sphere are held securely in place by internal male fasteners positioned on either or both hemispheres that engage corresponding female receptacles located on the opposing hemisphere.\nAlternatively, external fasteners are capable of securing the hemispheres together. Each hemisphere would have one or more screw portals that retain the screw head which will be described below in more detail.\nA bird must generate sufficient internally or externally applied forces to the two sections of the spherical avian cervical restraint collar simultaneously to distort the sphere to the point of disengaging the fastener from the receptacle of the locking mechanism. Removal of the collar by the bird in this manner is thus hindered by the overall properties and design of the sphere.\nClear, opaque, or colored semi-rigid non toxic materials may be utilized in construction of a sphere. Optional internal surface raised ribbing may be added to improve structural strength while maintaining a light weight characteristic.\nAs with prior art forms, different sizes of the spherical avian cervical restraint collar will be required to meet the physical variation of different sized birds. One size does not fit all. The avian cervical restraint collar requires selection of an appropriate size to fit an individual bird. As a general rule of thumb, the diameter of the sphere will approximate the widest diameter of the bird\"\"s head excluding the beak. It is further recognized that birds other than psittacines will not necessarily benefit from the application of one or more spherical restraint collars."}
{"text": "Prior to filling this application the most pertinent prior art known to me was:\nL. T. McBEE, U.S. Pat. No. 2,474,200;\nA. R. MOORE, U.S. Pat. No. 2,801,630;\nA. TEUFEL, U.S. Pat. No. 2,807,260;\nN. J. HALL, U.S. Pat. No. 2,820,455;\nA. CALABRESE, et al, U.S. Pat. No. 3,756,226."}
{"text": "The invention relates to a dispenser or an assembly suitable as a dispenser, serving as a receptacle, reservoir and/or discharger for media which may be liquid, pasty, powdery and/or gaseous. All components of the dispenser or assembly may be made of plastics or as compression or injection molded components. For discharge the dispenser can be freely held and simultaneously actuated single-handedly. Its length thus amounts to max. 10 cm or 7 cm, its largest width not more than 8 cm or 5 cm. The dispenser is suitable for dispensing single droplets of the medium, a jet or atomized particle or droplet aerosol thereof. Furthermore, the dispenser may be configured for discharging but a single dose of the medium or for a single stroke with no return stroke or for repeated discharges each with a spring-actuated return stroke inbetween.\nExperience has shown it to be expedient to compose complicated assemblies of components molded separately which during molding are located with or without a direct joint spaced away from each other or in a position other than that required in the operating condition. Reference is made to the German laid-open document 196 05 153 as well as to the pending German patent 198 13 078.3 in including the features and effects described therein in the present invention.\nThe invention is based on the object of providing a dispenser or a method of producing an assembly for a dispenser or the like which avoids the disadvantages of known configurations. It is more particularly the object to provide assemblies which have an increasing or decreasing inner or outer cross-section in the opposing direction. The dispenser is intended for facilitated production and safe operation.\nIn accordance with the invention two or more components are produced at the same time or with the same flow of plasticated material, immediately demolded once solidified or released in some other way at their jointing zones and then directly positioned relative to each other so that they can then be combined into an assembly. For the components the same material or differing materials may be employed. The components pass through the same temperature curves at the same time up to solidification and may have the same or differing volume of material. Expediently the components are produced in the same mold or so that they adjoin one or more common parts of the mold each integrally. This applies more particularly to the jointing surface areas of the components moldable juxtaposed in common by a movable part of the mold. After the components have solidified and subsequent retraction thereof or of another part of the mold these jointing zones are located exposed. The components can then be moved relative to each other until joined together and demolded completely where necessary. It is good practice when the components are located in production axially parallel or directly juxtaposed almost in contact with each other. Once the one component has been joined to the other it forms an elongation of the other component in the direction of its greatest extent. After being joined, forming the operating condition of the assembly for operation of the dispenser, the two components merge into a length which is smaller than the length of the one or other component. The components may, however, also be face joined without any mutual longitudinal engagement and locked in place mutually by a further component. Thus mutual locking of the components may be with zero clearance or positive, namely by being radially centered or by a captive lock.\nAlthough the configuration in accordance with the invention is suitable for the outer or base bodies of dispensers it is particularly expedient for core bodies. One such core body is located totally concealed in the interior of the dispenser of the corresponding base body, e.g. within a discharge nozzle. This base body may also form the third component for the cited locational lock. Advantageously one or both components of the assembly forms longitudinally a middle section of largest outer width, w adjoining at each end thereof an end section of comparitively reduced outer width. Each of the end sections is formed by another component. An end section may be a hollow needle having a smallest diameter at the tip of the needle of less than one millimeter and a length of less than 10 or 8 mm. The other end section may be a dished, fluted or outer face-recessed body having a radially protruding collar forming the shorter longitudinal part of the middle section.\nThe two components are advantageously joined to each other via a single connecting member or link directly joining each of the components by a link section. The link sections are then mutually movable and adjoined by a connecting location which may remain stationary in mutual movement of the components such as in movement of the corresponding link section relative to the corresponding component. The connecting location is expediently a hinging zone having a sole hinging axis and/or a designed frangible location at which the link sections are parted in mutual movement of the components and prior to attaining the operating position in forming opposing fractured surface areas. The mold cavity for the link may form the one or sole flow channel via which the plasticated material flows from the mold cavity for the one component, more particularly the larger volume component, into the mold cavity for the other component. The smallest cross-section of this channel and thus of the link may be less than 5, 2 or one tenth of a mm2.\nIn production the jointing surface areas of the two components later to directly adjoin in the operating position are expediently located in the same plane. Beyond one of these jointing surface areas a locking member or the like may protrude. In production these jointing surface areas may point in the same direction or in opposite directions. Up to each jointing surface area the link may also extend which may comprise a surface directly translating into the jointing surface areas in the same plane or frangible or parting surface areas in this plane after parting. The components may also be translated by a radial or linear movement into their operating position, the one component forming a sliding guide for the other component flanked only at the bottom and sides which, however, does not attain the guide until after a first portion of the shifting travel or after the link has been parted. In addition, the components may be produced separately and then assembled in accordance with the invention.\nIrrespective of the configuration as described, the dispenser is configured more particularly as a receptacle and reservoir for biological active substances over several weeks, months or even years. These may be physiological active substances containing hormones and/or cleavage products such as peptides containing protein. Such biological information transmitters which may contain amino acids and other similar active substances may be highly sensitive to moisture, this being the reason why they are held in the dispenser in a pressure-tight chamber which is not opened until immediately prior to delivery from the dispenser, e.g. by a closure being ruptured by means of the cited assembly.\nThese and further features of the invention also read from the description and the drawings, each of the individual features being achieved by themselves or severally in the form of sub-combinations in one embodiment of the invention and in other fields and may represent advantageous aspects as well as being patentable in their own right, for which protection is sought in the present."}
{"text": "Hydraulic type power steering has been known for a long time in which a hydraulic pump is driven by a belt on the shaft of the vehicle engine.\nHowever, that type of power steering suffers from numerous drawbacks.\nNumerous failures are to be observed in practice that are due to belt wear.\nIn addition, such a device occupies considerable axial space on the engine shaft.\nIt also requires the hydraulic pump to be located in an environment in which it is liable to be subjected to high levels of thermal and vibrational stress.\nElectric type power steering is also known in which an electric motor meshes with the steering column.\nNevertheless, such systems require complex electronics to control the electric motor.\nAlso, such systems are not of satisfactory reliability: the rotor of the electric motor can become jammed, thereby also jamming the steering column of the vehicle.\nTo mitigate those drawbacks, proposals have already been made for hydraulically assisted steering in which the pump is driven by an electric motor.\nNevertheless, assisted steering of the above type that has been known in the past makes use of an electric motor that is overdimensioned and some such steering systems require the electric motor to be powered continuously.\nSuch systems also suffer from a particularly high level of power consumption (500 kW to 1.5 kW).\nThis gives rise to a problem of dumping heat losses generated in the motor.\nTo solve that problem, proposals have already been made to control rotation of the electric motor as a function of the force imparted to the steering column. In this respect, reference can be made to patent application EP A 0 741 068.\nNevertheless, such a solution requires a torque pickup (such as a torsion bar) or a rotary speed pickup to be included on the shaft of the steering column."}
{"text": "In current technologies, the threshold voltage of semiconductor devices does not scale with the power supply voltage and ground rules because of the non-scalability of the sub-threshold slope. Thus, the minimum gate oxide thickness and/or maximum wordline boost voltage of the array MOSFET is constrained by reliability considerations.\nWhen used for the support MOSFET, the relatively thick gate oxide (having a thickness of greater than ≈6 nm for deep sub-μm technology) required by the array MOSFET results in degradation in the performance of the MOSFET device. Furthermore, if a thinner gate oxide is used to improve the performance of the support circuitry, charge transfer efficiency in the device array is compromised as a result of the reliability limitation of the wordline boost voltage.\nIdeally, in such technology, a dual gate oxide thickness is desired. In the prior art, it is known to subject the array transistor to a dual gate oxidation process or an alternative gate oxidation process as compared to the support circuitry. These additional gate oxidation processing steps are costly, and they are also yield limiting since one must utilize additional processing steps such as, but not limited to: masking, exposure, etching, oxidizing and strip masking, which grow a second oxide on the entire structure of the MOSFET device. As such, prior art gate oxidation processes are not reliable nor cost efficient.\nIn view of the drawbacks mentioned above with prior art processes of fabricating MOSFETs, there is a continued need for providing a new and improved method of fabricating a MOSFET and other devices in which a dielectric layer, e.g., gate oxide, having a dual thickness can be formed without adding extra processing steps and costs to the overall manufacturing process."}
{"text": "The present invention relates to a digital broadcast receiving apparatus that receives data on an electronic program guide (EPG) for explaining the contents of TV programs with at least character data and displaying the EPG and a method of displaying video data associated the TV programs in the EPG.\nThere are several known methods of displaying EPGs. One is to display only character data associated with TV programs. Another is to display program icons instead of character data so that even children who do not understand the character data can easily select TV programs. Still another is to display video data associated with TV programs which are often selected by a user so that he or she can easily discriminate his or her favorite TV programs from other programs in an EPG.\nThere are, however, several disadvantages to the EPG displaying methods described above.\nA known method of displaying only character data associated with TV programs tends to have much character data to explain the contents of the TV programs, so that users cannot understand the contents easily. A known method of displaying program icons instead of character data could offer child-friendly EPGs which, however, has a limit to explain the contents with icons, so that adults may have a difficulty in understanding the icons. A known method of displaying video data associated with user-favorite TV programs just offers images in an EPG irrespective of the size of the EPG or characters displayed with the images, thus lowering viewability.\nMoreover, video data has a larger amount of data than character data. Therefore, as the number of TV programs with which the video data is associated increases, the storage capacity of a memory that stores EPG data increases. In addition, when an EPG is scrolled, even though the EPG has no enough space for displaying video data and the video data cannot be displayed anymore after scrolled, the video data remains in the memory, which is a waste of memory."}
{"text": "A variety of different methods have been developed to assay oligonucleotides, including DNA or RNA fragments. Such assays are typically directed to determining whether a sample includes oligonucleotides having a particular target oligonucleotide sequence. In some instances, oligonucleotide sequences differ by only a few nucleotides, as in the case of many allelic sequences. Single nucleotide polymorphisms (SNPs) refer to alleles that differ by a single nucleotide. Even this single nucleotide difference can, at least in some instances, change the associated genetic response or traits. Accordingly, to determine which allele is present in a sample, the assay technique must be sufficiently sensitive to distinguish between closely related sequences.\nMany assay techniques include multiple components, each of which hybridizes to other component(s) in the assay. Non-specific hybridization between components (i.e., the hybridization of two non-complementary sequences) produces background noise in the assay. For example, closely related, but not identical, sequences can form imperfect duplexes in which base pairing is interrupted at positions where the two single strands are not complementary. Non-specific hybridization increases when the hybridizing components have similar sequences, as would be the case, for example, for many alleles and particularly for SNP alleles. Thus, for example, hybridization assays to determine which allele is present in a sample would benefit from methods that reduce non-specific hybridization or reduce the impact of non-specific hybridization on the assay."}
{"text": "It is known that certain electromagnetic oscillatory phenomena occurring in connection with meteorological conditions are correlated with variations of biological and pathological parameters.\nCharacterization of the respective pulses (Very-Low-Frequency Atmospherics; hereinafter: atmospherics), particularly according to frequency and wave form, has been undertaken, among others, by H. Baumer and J. Eichmeier in: Eine Anlage zur Registrierung der Atmospherics bei 10 und 27 kHz, Archiv fur Meteorologie, Geophysik und Bioklimatologie [a System for Recording Atmospherics at 10 and 27 KHz, Archieves for Meterology Geophysics and Bioclimatology], Ser. A, 29, 143 to 155 (1980), and by W. Sonning, H. Baumer and J. Eichmeier in: Die Atmospherics-Aktivitat bei 10 und 27 kHz als Indikator fur die Dynamik der tropospharischen Wettervorgange [Atmospherics Activity at 10 and 27 KHz as an indicator for the Dynamics of Tropospheric Weather Processes], Archives For Meteorology, Geophysics and Bioclimatology, Ser. B, 29, pp. 299-312 (1981). Atmospherics are short-time pulses in the form of damped oscillations. Their respective duration comprises approximately 5 to 6 half-waves. They are a mixture of different frequencies. Using measuring techniques, atmospherics have been determined by Baumer according to their predominant frequency fraction of e.g. 6, 8, 10, 12 and 28 kHz. They will therefore be designated in the following as \"according to Baumer\"=atB.\nA relationship has already been established between the frequency of occurrence of some definite atmospherics and various biological parameters. H. Baumer, in: Die Meteotropie eines Dichromat-Gelatinesystems, Technischer Informationsdienst des Bundesverbandes Druck e. V. [The Meteotropy of a Dichromatic-Gelatine System, published by the Technical Information Service of the Federal Association for Printing, Registered Assocation] II/1982, pp. 1 to 17, describes the difference in pulse rates in the 10 kHz and 28 kHz ranges and how they distinctly correlate with the diffusion behaviour of gelatine, which, in turn, influences the production process of rotogravure forms; cf. also H. Baumer and J. Eichmeier: Relationship between the Pulse Rate of Impulsstrahlung and the Diffusion Time of Ions in Gelatine Films, Int. J. Biometeor. 1980, Vol. 24, No. 3, pp. 271 to 276, as well as H. Baumer and J. Eichmeier: The Biophysically Active Wave Forms of Impulsstrahlung Incident on Gelatine Films, Int. J. Biometeor. 1982, Vol. 26, pp. 85-90.\nThe correlation of atmospherics with certain diseases such as epilepsy and cardiac infarction, for which some relationship with meteorological conditions had already been supposed before, has been described in the EP 0 120 991 A2 (U.S. Pat. No. 4,631,957); cf. also G. Ruhenstroth-Bauer, H. Baumer et al.: Epilepsy and Weather: A significant Correlation Between the Onset of Epileptic Seisures and Specific Impulsstrahlung - A Pilot Study, Int. J. Biometeor. 1984, Vol. 28, No. 4, pp. 333-340, and G. Ruhenstroth-Bauer, H. Baumer, among others: Myocardial Infarction and the Weather; A significant Positive Correlation between the Onset of Heart Infarct and 28 kHz Impulsstrahlung - A Pilot Study, Clin. Cardiol., 8, pp. 149-151 (1985).\nFurthermore, a correlation has been found to exist between the 8 kHz and 10 kHz atmospherics and inflammatory processes in rats; see G. Ruhenstroth-Bauer, O. Rosing, and H. Baumer, Naturwissenschaften (Natural Sciences) 73, p. 625 (1986).\nIn addition, an important correlation between natural atmospheric spectra and the in vitro incorporation of (.sup.3 H)-thymidine into the nuclear DNA of C6-glioma-cells has been stated; cf. Vogl, G. Hoffmann, B. Stopfel, H. Baumer, O. Kemsky and G. Ruhenstroth-Bauer: Significant Correlations Between Atmospheric Spectra According to Baumer and the in vitro Incorporation of (.sup.3 H)-Thymidine into the Nuclear DNA of C6-Glioma-Cells, FEBS Letters, Vol. 288, No. 1, 2, pp. 244-246 (1991).\nA summary of the relations between atmospherics and biological and pathological parameters known so far is to be found in G. Hoffmann, S. Vogl, H. Baumer, O. Kemsky and G. Ruhenstroth-Bauer: Significant Correlations Between Certain Spectra of Atmospherics and Different Biological and Pathological Parameters, in: Int. J. Biometeorol. (1991) 34, pp. 247-250.\nIt is not possible as yet, however, to detect any relation of causality between the occurrence of those atmospherics and the described biological and pathological parameters."}
{"text": "This invention relates to the detection and signal processing of electrical brain activity.\nThe purposes of this invention include the detection of information processing undertaken in the brain, the detection of concealed information in the brain, communication from the brain to a computer, and command and control of computers and electronic and mechanical equipment by the brain.\nThe Farwell MERA System is a new technology for the detection of concealed information that revolves around the non-invasive recording of brain electrical activity. The electrical brain activity pattern recorded and of interest is a specific multifaceted electroencephalographic response (MER) that occurs immediately after an examinee is visually presented (via a computer screen) with words, short phrases, acronyms, or pictures that are recognized and cognitively processed by that subject. This phenomenon, coupled with its absence following the presentation of the same information to a subject for whom the material is unknown or irrelevant, is the basis for discriminating between subject guilt and innocence. This would potentially allow for the determination of a whole host of issues of interest to the law enforcement and intelligence communities, e.g., (1) does a suspect have guilty knowledge connecting him to specific investigated criminal activity, (2) does an intelligence source have knowledge of the internal workings of a hostile intelligence agency that would indicate that he was an intelligence officer of that agency and not who he claimed to be, (3) has an informant, a debriefed spy, or a suspected member of a criminal organization accurately described the entirety of his actions and knowledge, (4) did a convicted serial killer who claims to have killed 40 to 50 individuals, other than the one(s) he was convicted of, actually commit these acts, or are these claims merely the bravado of a condemned prisoner.\nThe potential benefit of this program extends to a broad range of law enforcement applications, including organized crime, violent crime, white-collar crime, drug-related crime, foreign counterintelligence, non-traditional targets, and other categories of casework as well. This new technology promises to be of tremendous benefit both at the national level and for state and local law enforcement agencies.\nThis application describes a technology that is capable of detecting concealed information stored in the brain through the electrophysiological manifestations of information-processing brain activity. Additional information is described in a previous patent application of the inventor, U.S. patent application Ser. No. 08/016,215, entitled \"Method and Apparatus for Truth Detection\" filed on Feb. 11, 1993, which is expressly incorporated herein by reference.\nThis technique provides a means of distinguishing guilty and innocent individuals in a wide variety of law enforcement and information detection situations. The research described below demonstrates that the system is also effective in distinguishing between members of a particular organization (in this case, the FBI) and others who are not knowledgeable regarding that organization.\nWhen a crime is committed, traces of the event are left at the scene of the crime and elsewhere. The task of the investigators is to reconstruct what has happened and who has been involved, based on the collection of such evidence.\nIn addition to the physical and circumstantial evidence that can be obtained, there is one place where an extensive record of the crime is stored: in the brain of the perpetrator. If this record could be tapped, criminal investigation and counterintelligence could be revolutionized.\nUntil recently, the only method of attempting to discern what information regarding a crime or other situation of interest was stored in the brain of a suspect or witness has been (1) to interrogate the subject, and (2) to attempt to determine whether or not the subject is lying.\nConventional control question (CQT) polygraphy has been used as an aid in the attempt to detect deception in such reports. The fundamental theory of conventional polygraphy is that a deceptive individual will be more concerned with and experience more emotional arousal in response to relevant questions than control questions, and this emotional arousal will be accompanied by corresponding physiological arousal which can be measured. Traditional interrogative polygraph (\"lie detection\") methods rely upon using questioning formats in conjunction with the recording of physiological parameters that reflect autonomic nervous system (ANS) activity (e.g. blood pressure, heart rate, sweating, etc.). This information is peripheral to the cognitive aspects of deception or of concealing guilty information."}
{"text": "1. Field of the Invention\nThe present invention relates to accessing to information over the Internet. In particular, the present invention relates to a customized access to information over the Internet by various internet appliances with various processing capabilities.\n2. Discussion of the Related Art\nAs the Internet has become a preferred medium for information access and dissemination, many different devices (e.g., mobile phones, personal digital assistants and handheld computers) can now be used to access information on the Internet. In general, these devices typically have much lesser text and graphical processing capabilities than a conventional desktop computer. (For convenience, in the remainder of this description, these devices are collectively referred to as “internet appliances”.) As much of the information on the Internet is organized for access by a desktop computer using a hypertext protocol (e.g., http), access to such information by a device other than a desktop computer can be inefficient. For example, many web pages are designed with a high-resolution graphical display in mind. Even when possible, accessing such web pages from a mobile telephone without a graphical display and providing only a limited number of short lines for text display can be a very frustrating experience.\nTo accommodate the different capabilities of the internet appliances, in the prior art, an operator of a website typically provides for each supported internet appliance a specialized “edition” of the website accessible through a specialized gateway. For example, since the current generation of mobile telephones are typically only capable of displaying text of a small number of characters per line, an operator would provide specially designed text-only “stripped down” web pages accessible through a wireless access protocol (WAP) gateway. In most instances, information available in the general edition of the web pages are included or excluded by the designer or operator based on its resource availability or other criteria, without user participation. Often, therefore, information important to some users is arbitrarily excluded, thereby severely reducing the utility of the web pages.\nWhere a specialized website is not available, the gateway would provide only the text from the web pages and discard or ignore graphical information, animation or other functions embedded in the web pages. In such an instance, no attempt is typically made to filter the information based on the content of a web page. Consequently, a relatively small web page can result in the user pressing the “scroll” key a large number of times. Many users therefore do not consider internet appliances to be suitable for serious information retrieval purposes."}
{"text": "1. Field of the Invention\nThe present invention is generally in the field of electrical circuits and systems. More specifically, the present invention is in the field of memory devices.\n2. Background Art\nIn various memory devices, such as RAMs (random access memory), CAMs (content addressable memory), or other memory devices, that handle a large amount of input data, the number of inputs to the memory device need be managed and reduced to, for example, reduce problems associated with data routing when there are numerous input pins. However, the reduction of the number of inputs need be accomplished without impairing the ability of the memory device to handle a large amount of input data. Moreover, it is desirable to preserve flexibility so that various system design requirements in receiving input data can be met.\nFor example, if a system is capable of providing write or compare data at a high rate into a RAM or a CAM, then it is desirable to utilize the high data rate advantageously so that data can be received by the RAM or CAM at a faster pace. On the other hand, if a system is not capable of providing data at a high rate into a RAM or CAM, then it is desirable to accommodate the slower data rate as well. In this manner a system designer can utilize fewer input pins of the RAM or CAM for providing input data, or can use a more conventional data rate which requires a greater number of input pins. The capability to provide and receive data at higher rates results in speed improvement, and a reduction of input pins and reduces input/output (I/O) routing complexity in memory devices, such as RAMs and CAMs.\nThus, it is desirable to maintain flexibility to accommodate receiving data at various data rates, and to switch between a slower incoming data rate and a faster incoming data rate in various memory devices, such as RAMs, CAM, as well as other memory devices."}
{"text": "In recent years, in a DC feeding system, regenerative power generated by a regenerative brake of a train is used as power-running electric power for other trains via an overhead wire. In such a DC feeding system, surplus regenerative power is intermittently generated in a same power transformation zone when the regenerative power exceeds the power-running electric power and it is effectively re-used by a power regenerative inverter installed in a substation.\nMeanwhile, such a technique is disclosed in which if supply of commercial frequency power to an AC bus bar of the substation for electric railways is shut off, the operation of the power regenerative inverter is stopped and the power regenerative inverter is operated as a self-exciting inverter; and power supplied from an adjacent DC substation via the DC feeding system is converted to AC power. Accordingly, the emergency power is supplied to station building facilities of a station building via a high-voltage or extra high-voltage distribution system (Patent Literature 1, for example)."}
{"text": "A. Field of Invention\nThis invention pertains to a method and apparatus for operating an internal combustion engine using a fuel consisting of water and a water-soluble flammable substance that is injected into a mixture of hydrogen and air.\nB. Description of the Prior Art\nThe use of fossil fuels to run engines that used, for example, in cars and other vehicles, as well as many other engines used for a variety of purposes, is based on a very old concept based on the internal combustion engines developed in the nineteenth century. Despite intense research and development for alternate fuels for the last 50 years, fossil fuel derived from petroleum or natural gas, is still essentially the primary source of energy almost all the internal combustion engines presently in use all over the world.\nAs a result, the world supply of fossil fuels have been severely depleted creating a shortage, and the price of oil has been climbing for the past 40 years. In addition such fuels are very polluting and some suggest that it has either been the primary cause or has contributed substantially to global warming. All these factors led to many efforts to find and harness renewable energy sources other than traditional fossil fuels. Several alternative fuels have been introduced in the past few years to reduce the impact of petroleum depletion, including hybrid cars, electric cars, bio diesel, hydrogen based cars, etc. However, none of these solutions were effective. One reason for this lack of success is that they require a completely new infrastructure for the production of the engines, as well as the production and distribution of the fuel. Moreover, the most solutions proposed so far were incompatible with the existing engines and, therefore. The cost of replacing all the existing fossil burning engines may be so high that it may render any solution based on alternate fuels unacceptable, at least, in a short term basis.\nWater as a source of fuel has been suggested by many in the past and many experiments have been conducted testing such systems. The basis of such experiments is the fact that water can be separated in to hydrogen and oxygen and the resulting stoichiometric mixture can be fed in to an internal combustion engine to generate power. However past experiments yielded unsatisfactory results. The main obstacle for their success is based on the fact that the energy required to separate the water into its components is much greater than the energy produce by the engine. In addition the amount H2 mixture needed to run a typical automotive engine is too large to make such a system practical.\nSystems are presently available on market that can be used as accessories or add-ons to internal combustion engines using fossil fuels, however independent tests have shown that, in fact, these systems have very little, if any, effect on the overall efficiency of the engine.\nA system developed by the present inventors is described in two co-pending applications includes means of generating from water and supplying a small amount of hydrogen/oxygen gas mixture into a standard internal combustion engine. (See U.S. Patent Application Publications 2010/0122902 and 20110203917). More specifically, these co-pending applications describe an efficient process and apparatus for generating a two-to-one mixture of hydrogen and oxygen, commonly referred to a brown gas or HHO. The mixture helps increase the efficiency of the conventional internal combustion engine by burning the fossil fuel more efficiently. While this latter system is much more efficient that previously described systems; its efficiency is still limited by the amount of hydrogen and oxygen produced on board a vehicle. Moreover, the internal combustion engine described is still burning a fossil fuel."}
{"text": "With the popularity of digital photography and digital image processing, consumers have increasingly desired to transfer photographic images stored on conventional film negatives into electronically stored digital images. Typically, this is accomplished by loading a sheet of processed film into a scanner and scanning the film to produce the digital image. Processed film is normally cut into sheets containing one to six images. Thus, if a user has a large number of negatives to scan, the process of loading each individual sheet of film into the scanner can become overly time-consuming. Accordingly, there is a desire for improved systems and methods for automate the loading and scanning of multiple sheets of film.\nConventional systems for handling the feeding of paper or film documents, such as those used in photocopiers, printing presses, printers and scanners, are not well suited for the handling of film. In particular, the rollers used for feeding individual sheets from a stack of paper or film may damage the image on sheets of film. In addition, these loading mechanisms are configured to load a large number of identically-sized sheets of paper in standard sizes such as 8.5″×11″ or 8.5″×14″. In contrast, photographic film negatives are often manually cut, resulting in sheets of film of varying lengths that are difficult to accurately load on a bulk basis. In addition, photographic film can change its shape over time or during operation, such as when the film curls around unpredictable angles."}
{"text": "An electronic device housing device that can house a plurality of electronic devices, such as a hard disk drive (HDD), a power source device, and a circuit substrate, inside a housing is known in the art. Examples of the electronic device housing device include a storage device equipped with a plurality of HDDs, a control device that controls the HDDs, and a power source device.\nIn recent years, there has been an increasing demand for thin storage devices that can house several HDDs arranged side by side in a row on one to three shelves. The thin storage devices are mainly used as placed on a rack. Such storage devices are called “rack-mounted storage devices”.\nA plurality of slots partitioned by a plurality of partition plates are provided in the housing of the storage device to allow a cartridge HDD carried on a carrier with a grip to be housed in each of the slots.\nAn HDD is a storage device that magnetically stores data utilizing a rotatable magnetic disk medium and a magnetic head. Mechanical parts such as the magnetic disk medium and the magnetic head are housed in a sealed case.\nA connection substrate formed by a single plate and equipped with a plurality of connection terminals is disposed near the center of the housing of the storage device. The plurality of connection terminals are disposed in parallel on the connection substrate to face the slots, and electrically connected to external connection terminals and power source terminals of the HDDs.\nThe connection substrate is provided near the center portion of the housing as discussed earlier, and therefore divides the housing space inside the housing into a front housing portion and a rear housing portion. The plurality of slots are disposed in the front housing portion, and the control device and the power source device electrically connected to the connection substrate are disposed in the rear housing portion. Each of the HDDs sends and receives a signal to and from the control device, and is supplied with power from the power source device, via the connection substrate.\nA cooling fan is provided in the rear housing portion of the housing. The cooling fan is driven to allow a gas to be sucked into the slots from gas suction ports provided in a front panel of the housing and the carriers.\nThe sucked gas flows straight into the rear housing portion while cooling the HDDs in the slots, cools the control device and the power source device, and thereafter is discharged from an exhaust port provided in the rear wall surface of the housing. The HDDs, the control device, and the power source device are cooled by such an air-cooling structure. However, the plate-shaped connection substrate provided near the center portion obstructs the gas flow linearly guided from the gas suction ports to the rear housing portion across the connection substrate.\nExamples of the related art described above are disclosed in Japanese Laid-open Patent Publication No. 2009-170649 and Japanese Laid-open Patent Publication No. 11-204974.\nIn recent years, it has been further desired to reduce the thickness of a storage device. Reducing the thickness, however, makes it difficult to secure a passage of a gas for cooling electronic devices."}
{"text": "This invention relates to a competitive dice game which is Played with elements which are portable so that the game can be played anywhere and provides a unique and exciting competitive game experience.\nThere exist numerous games employing dice and other elements, frequently using game boards and pieces for moving the pieces around the game boards. One popular game is backgammon which is played with dice, a doubling cube, a board, pieces which move on the board and other associated elements. This game requires several independent elements and is difficult to easily transport and use without inconvenience.\nThere are numerous dice games in which dice are rolled, especially in betting environments. These games involve a single roller or player who is seeking to achieve a specific number with each roll, with there being no cumulative roll or count as the dice are repeatedly rolled by individual players in the game.\nAn object of this invention is to provide an improved competitive dice game which is susceptible to widespread use, easily portable and full of excitement and strategy suitable to players of different skill levels."}
{"text": "In the cathodic protection of ferrous structures, especially pipelines, the use of a mixture of alkali bentonite, gypsum and sodium sulfate as a backfill for underground magnesium-base anodes is well known, the particulars of which are shown in the patents listed below. It is noted that among the teachings in the patents it is taught that \"alkali bentonite\" is the operable form of bentonite, but that \"alkaline earth bentonite\" is inoperable.\nU.S. Pat. No. 2,478,479 discloses a magnesium-base alloy on a Mg-Al alloy core, buried in a backfill of bentonite-gypsum mixture, for galvanic protection of a ferrous metal pipeline.\nU.S. Pat. No. 2,480,087 discloses a backfill consisting of naturally-occurring \"bentonite\" in admixture with gypsum and a water-soluble metal salt, such as sodium sulfate. The operable bentonite is said to be \"alkali bentonite\" in contradistinction to \"alkaline earth bentonite\" which is said to be inoperable.\nU.S. Pat. No. 2,525,665 discloses a gypsum-bentonite-sodium sulfate backfill such as is described in U.S. Pat. No. 2,480,087 above.\nU.S. Pat. No. 2,527,361 discloses a gypsum-bentonite-sodium sulfate backfill such as is described in U.S. Pat. No. 2,480,087 above.\nU.S. Pat. No. 2,567,855 discloses a backfill of gypsum-bentonite-sodium sulfate.\nU.S. Pat. No. 2,601,214 discloses a backfill comprising a major proportion of magnesium sulfite and a minor proportion of \"sodium-type\" bentonite (montmorillonite).\nA reference for mineralogical information about bentonite clays, and other clays of the montmorillonite type, is \"Applied Clay Mineralogy\" by Ralph E. Grim, published by McGraw-Hill Book Company, Inc., New York, 1962.\nAs used in this application the term \"bentonite\" is used in referring to minerals which are largely composed of montmorrillonite clays such as are mined as alterations of volcanic ash, and the like. Alkali metal bentonites (e.g., sodium bentonite) are known to swell upon addition of water, and to contract or de-swell upon removal of water, in contradistinction to alkaline earth metal bentonites (e.g., calcium bentonite) which undergo little, if any, such swelling or de-swelling."}
{"text": "Wireless communication devices have become smaller and more powerful in order to meet consumer needs and to improve portability and convenience. Consumers have become dependent upon wireless communication devices such as cellular telephones, personal digital assistants (PDAs), laptop computers, and the like. Consumers have come to expect reliable service, expanded areas of coverage, and increased functionality. Wireless communication devices may be referred to as mobile stations, stations, access terminals, user terminals, terminals, subscriber units, user equipment, etc.\nA wireless communication system may simultaneously support communication for multiple wireless communication devices. A wireless communication device may communicate with one or more base stations (which may alternatively be referred to as access points, Node Bs, etc.) via transmissions on the uplink and the downlink. The uplink (or reverse link) refers to the communication link from the wireless communication devices to the base stations, and the downlink (or forward link) refers to the communication link from the base stations to the wireless communication devices.\nWireless communication systems may be multiple-access systems capable of supporting communication with multiple users by sharing the available system resources (e.g., bandwidth and transmit power). Examples of such multiple-access systems include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, and orthogonal frequency division multiple access (OFDMA) systems."}
{"text": "The invention of this application relates to an improvement in the design of rack in which vertical beams and horizontal columns are attached together by bolts, for simple assembly and disassembly. Welded racks may be strong, but they cannot be easily disassembled for moving, and the labor involved in welding them together is substantial.\nBolted racks, while easily assembled, exhibit the problem that any looseness in the beam to column connection will result in a leaning, wobbling rack. This in turn induces higher stresses in many of the components which might cause collapse of a heavily loaded rack. Also, problems may arise because of irregularities in the floor of the warehouse or other place where the rack is positioned. This, along with small errors in the length and width of the beams and columns as well as the positioning of bolt holes therein, can result in the creation of unpredictable positional variations in the beams and columns as one attempts to affix them together. This, in turn, may create significant difficulties in getting a tightly bolted connection between the beams and columns, since the parts may not quite fit precisely together because of the unpredictable positional variations.\nThus, the final result may be a less than satisfactory leaning or wobbling rack, which may be quite unacceptable in any circumstances.\nIn accordance with this invention means are provided for affixing beams and columns together by bolt means in a tight, rigid manner despite unpredictable positional variations in the beams and columns as they are so affixed.\nWhile racks are specifically contemplated as the environment for the use of this invention, other structures besides racks may take advantage of this invention as well, such as frames, buildings, mezzanines, or any other structure using a beam to column connection, or any connection between structural members."}
{"text": "1. Technical Field\nThis invention relates generally to joints for linking relatively movable vehicle steering components to one another, such as ball joints, tie rod ends, and sway bar links.\n2. Related Art\nVehicle suspension systems and steering systems typically include joints, such as tie rod end ball-type joints for operable attachment of a tie rod end to a steering knuckle and a ball joint for coupling the steering knuckle to a control arm. In addition, other applications, such as carnival rides or any other mechanism with relatively movable joints, typically have ball joints to facilitate the relative movement between linked components. Such ball joints typically include one or more bearings that are received in a housing and a ball stud that slidably contacts the bearing or bearings to allow the housing and ball stud to articulate relative to one another during use.\nUpon assembly of ball joints, it is generally desirable to build in a frictional resistance between the ball stud and housing that is within a predetermined torque tolerance. In addition, it is essential that the ball joints exhibit a long and useful life, and of additional importance, it is important that the ball joints be economical in manufacture. If the frictional resistance or torque is too high, it may impede the motion of the mechanism and/or make installation difficult. If the frictional resistance is too low, it may result in an undesirable “out-of-box feel” to which the installer of the joint will believe the socket to have excessive looseness, and therefore, shorter operating life.\nIt is known to construct ball joints from metal, including coated metal bearings against which the metal ball stud pivots. However, although the coated metal bearings can provide a desirable “out-of-box” feel and exhibit a long and useful life, they typically come at a high cost in manufacture.\nIn an effort to reduce costs associated with manufacture, it is known to construct tie rod end ball joints with glass-filled nylon or fiber-reinforced nylon bearings against which a metal ball stud pivots. Although the cost of manufacture is greatly reduced, the glass-filled nylon bearings provide a reduced useful life as compared to metal bearings.\nDuring manufacture of such socket assemblies, prior to molding a nylon bearing, the nylon resin must be sufficiently dried to remove any water from the nylon resin and minimize water in the socket assembly that could corrode the metal ball stud. Even with this drying operation, the metal ball studs of socket assemblies with nylon bearings are coated with a corrosion resistant material for further corrosion protection as water often inevitably finds its way into the socket assembly, if not during assembly, then during use."}
{"text": "Voice messaging systems that enable users to send and retrieve voice mail messages are known in the communication arts. In a typical prior art voice messaging system a telephone is connected to a private branch exchange (PBX) that may utilize a notification mechanism, such as a message waiting indicator light, to notify a message recipient that a new message is waiting for them. Many wireless telephone communication systems also provide a Short Message Services (SMS) feature that allows users to send and/or receive short text messages. Today, many modern communication systems provide messaging services via packet-based networks, i.e., those that operate in accordance with the Internet Protocol (IP). A unified messaging (UM) system handles voice, facsimile, regular text messages, and computer-readable documents as objects in a single mailbox that a user can access either with a regular email client, or by telephone. A UM system typically connects to a PBX to provide automated attendant, audiotext, and voice mail services to subscribers or users. For instance, a personal computer (PC) user with multimedia capabilities typically can open and playback voice messages, either as speech or text. Unified messaging is therefore particularly convenient for mobile business users because it allows them to reach colleagues and customers through a PC or telephone device, whichever happens to be available.\nThere are times when a caller connected to a messaging system via a voice-only channel leaves a voicemail, but also would like to attach media content such as electronic mail (email), Web pages, financial data, documents, and/or video attachments to the voicemail message. But due to the limitations of existing UM and telephony systems users are unable to attach additional rich media content (e.g., documents, media clips, Uniform Resource Locators (URLs), etc.) to a unified message that is left as a voicemail."}
{"text": "The present invention generally relates to an automatic implantable atrial defibrillator for delivering cardioverting or defibrillating electrical energy to the atria of a human heart. The present invention is more particularly directed to such an atrial defibrillator which inhibits therapy delivery when the time span between detected atrial events exceeds a predetermined limit to avoid delivering therapy when the heart reverts to normal sinus rhythm or during a normal sinus rhythm cycle of the heart.\nAtrial fibrillation is probably the most common cardiac arrhythmia. Although it is not usually a life threatening arrhythmia, it is associated with strokes thought to be caused by blood clots forming in areas of stagnant blood flow as a result of prolonged atrial fibrillation. In addition, patients afflicted with atrial fibrillation generally experience palpitations of the heart and may even experience dizziness or even loss of consciousness.\nAtrial fibrillation occurs suddenly and many times can only be corrected by a discharge of electrical energy to the heart through the skin of the patient by way of an external defibrillator of the type well known in the art. This treatment is commonly referred to as synchronized cardioversion and, as its name implies, involves applying electrical defibrillating energy to the heart in synchronism with a detected ventricular electrical activation (R wave) of the heart. The treatment is very painful and, unfortunately, most often only results in temporary relief for patients, lasting but a few weeks.\nDrugs are available for reducing the incidence of atrial fibrillation. However, these drugs have many side effects and many patients are resistant to them which greatly reduces their therapeutic effect.\nEarly implantable atrial defibrillators were proposed to provide patients suffering from occurrences of atrial fibrillation with relief. Unfortunately, to the detriment of such patients, these early atrial defibrillators never became a commercial reality.\nTwo such proposed defibrillators, although represented as being implantable, were not fully automatic, requiring human interaction for cardioverting or defibrillating the heart. Both of these proposed defibrillators required the patient to recognize the symptoms of atrial fibrillation. One defibrillator required a visit to a physician for activation of the defibrillator and the other defibrillator required the patient to activate the defibrillator from external to the patient's skin with a magnet.\nAn improved atrial defibrillator which provides automatic operation is fully disclosed in U.S. Pat. No. 5,282,837 entitled \"IMPROVED ATRIAL DEFIBRILLATOR AND METHOD,\" and which issued on Feb. 1, 1994 in the names of John M. Adams and Clifton A. Alferness. This patent is assigned to the assignee of the present invention and is incorporated herein by reference.\nIn addition to being automatic in operation, the atrial defibrillator of the above-referenced patent includes further features to assure the safe operation of the device. For example, to assure that the cardioverting electrical energy is not applied during the T wave vulnerable period of the heart, the atrial defibrillator provides R wave detection of increased reliability which is utilized to advantage in synchronizing the delivery of the cardioverting electrical energy to the atria with an R wave of the heart. Further, as another feature, an electrode system is utilized which minimizes the amount of energy that is required to cardiovert the atria. This is achieved by locating the cardioverting electrodes in or near the atria of the heart to provide a cardioverting energy path which confines substantially all of the cardioverting electrical energy to the atria of the heart.\nFurther improvements directed to the safe operation of an implantable automatic atrial defibrillator are described in U.S. Pat. No. 5,207,219 which issued on May 4, 1993 for \"ATRIAL DEFIBRILLATOR AND METHOD FOR PROVIDING INTERVAL TIMING PRIOR TO CARDIOVERSION,\" and which is also assigned to the assignee of the present invention and incorporated herein by reference. The atrial defibrillator there disclosed provides an answer to the observation that during episodes of atrial fibrillation, the cardiac rate increases to a high rate and/or becomes extremely variable. At high or variable cardiac rates, the R wave of a cardiac cycle may become closely spaced from the T wave of the immediately preceding cardiac cycle. This creates a condition known in the art as an \"R on T\" condition which is believed to contribute to induced ventricular fibrillation if the atria are cardioverted in synchronism with the R wave close to the preceding T wave. In order to prevent cardioversion of the atria during an R on T condition, the atrial defibrillator described in U.S. Pat. No. 5,207,219 detects for a cardiac interval longer than a minimum interval prior to delivering the cardioverting electrical energy to the atria. This assures that the cardioverting electrical energy is not delivered during an R on T condition.\nIn addition to the foregoing, there is certain electrogram data related to atrial fibrillation detection and cardioversion from which the cardiologist would benefit. Such information includes electrograms of the heart during fibrillation to confirm proper operation of the atrial fibrillation detector, electrograms of the heart prior to cardioversion and electrograms of the heart from immediately prior to and ending after the deliverance of the cardioverting electrical energy to the atria to confirm that the application of the cardioverting electrical energy was synchronized with an R wave and not on a T wave and to also confirm that the cardioversion was successful by returning the heart to normal sinus rhythm. One such implantable automatic atrial defibrillator capable of storing such electrogram data for later recall by the cardiologist is fully disclosed in U.S. Pat. No. 5,522,850 which issued on Jun. 4, 1996 for \"DEFIBRILLATION AND METHOD FOR CARDIOVERTING A HEART AND STORING RELATED ACTIVITY DATA,\" which is assigned to the assignee the present invention, and which is incorporated herein by reference.\nOne problem with storing electrogram data related to the cardioversion of atrial fibrillation is that during atrial fibrillation, the heart may spontaneously convert to normal sinus rhythm or do so for at least one or more cardiac cycles. Such a cardiac cycle would most likely satisfy electrical energy therapy criteria, including a minimum interval criteria. Hence, it is possible for the defibrillator to correctly detect the presence of atrial fibrillation and then deliver cardioverting energy on cycle which, when later reviewed by the cardiologist, would appear to be a normal sinus rhythm cardiac cycle."}
{"text": "1. Field of the Invention\nThe present invention relates to elements for the quantitative or semi-quantitative analysis of liquids which may contain gentisic acid as an interferent.\n2. Description of the Related Art\nIt is well known in the art to perform a quantitative or semi-quantitative analysis of a liquid by contacting that liquid with an analytical element containing reagents capable of yielding a detectable product in proportion to the concentration of a predetermined analyte in the liquid. One particularly useful method involves an enzymatic assay wherein the predetermined analyte, upon contact with the analytical element, is oxidized in the presence of an enzyme contained therein to produce a peroxide in proportion to the concentration of the predetermined analyte in the liquid undergoing analysis. A detectable product is then yielded by the reaction of the peroxide with an indicator composition in the presence of a substance having peroxidative activity. This detectable product should be formed in direct proportion to the peroxide present and thus also in proportion to the concentration of the predetermined analyte. Elements and analyses of this type are described in U.S. Pat. No. 3,992,158 and in a copending U.S. application by B. J. Bruschi, Ser. No. 712,972, filed Aug. 9, 1976, both of which are incorporated herein by reference.\nMethods of analysis employing reaction mechanisms other than the above-described peroxide mechanism to produce a detectable product are also known. For example, U.S. Pat. No. 3,711,252, describes a method for the quantitative analysis of uric acid in aqueous liquids wherein the aqueous liquid is contacted with a carrier element containing a ferric salt and either 2,4,6-tri(2-pyridyl)-1,3,5-triazine of 2,2':6',2\"-terpyride, in a buffered acidic medium. A color change is produced which is directly proportional to the concentration of uric acid in the aqueous liquid.\nAdditional methods are described in U.S. Pat. No. 3,801,466.\nIn all of the above-cited references to elements and methods for their use, it is also recognized that substances present in the liquid undergoing analysis other than the predetermined analyte may interfere with or bias the analytical reactions such that the detectable product is not formed in direct proportion to the predetermined analyte alone. This is particularly true for relatively low concentration analytes. For example, in analyses for uric acid or lactic acid in aqueous liquids such as serum or urine, it is recognized that gentisic acid can interfere with the reactions used to indicate the concentrations of uric acid or lactic acid. This is a significant problem, because it is well known that gentisic acid may often be present in liquids such as serum or urine.\nGentisic acid is a metabolic product of acetylsalicylic acid (aspirin) and would be expected to be found in the body fluids of anyone who has recently ingested common aspirin. Its presence has been recognized and methods are available for its analysis in liquids. Since uric acid and lactic acid are normally present in body fluids in relatively small concentrations, the recent ingestion of one or two doses of aspirin can effectively destroy the accuracy of a lactic acid or uric acid analysis.\nIn some of the analytical methods described above, such as those discussed in U.S. Pat. Nos. 3,711,252 and 3,801,466, gentisic acid is falsely detected as more of the predetermined analyte, because it reacts with the analytical reagents to form the detectable product in the same way that the predetermined analyte does. The result is a false indication that there is a higher concentration of predetermined analyte than is actually present.\nMethods are available and known to avoid interferences of this type. For example, U.S. Pat. No. 3,711,252, suggests prevention of gentisic acid interference by incorporation of persulfate in the analytical element. U.S. Pat. No. 3,801,466 suggests a multi-step method of avoidance involving preparation of comparative test samples in one of which the predetermined analyte is totally eliminated by a pre-analysis reaction. The two samples are then analyzed for predetermined analyte, and the difference in results between the two indicates the concentration of interferents such as gentisic acid that may be present. While such methods of avoidance are useful, they are either inconvenient to use (involving multiple steps) or are applicable only to one method of analysis. For example, the use of persulfate suggested by U.S. Pat. No. 3,711,252, would not be successful in avoiding gentisic acid interference with the proxide-linked analyses described previously.\nThe mechanism of interference of gentisic acid with a peroxide-mechanism-type analysis is quite different from the interferences described above, wherein gentisic acid is falsely detected as predetermined analyte. In a peroxide-mechanism-type analysis, gentisic acid interference produces the opposite effect. It causes a false indication that there is a lower concentration of predetermined analyte than is actually present. Unlike the other analytical methods wherein gentisic acid reacts to form the detectable product just as the predetermined analyte does, in the peroxide-linked analyses gentisic acid competes with the indicator composition in the presence of a substance having peroxidative activity in order to react with the peroxide formed by the interaction of predetermined analyte and enzyme. Thus, less peroxide is available to react with the indicator composition to produce the detectable product, and the concentration of predetermined analyte indicated is falsely low.\nAlthough the mechanism of the competition between the indicator composition and gentisic acid for peroxide is not definitely known, the following hypothesis is presented as a possible explanation of the interference.\nIn the peroxide-linked analyses which use a analytical element as a test-reagent carrier, all of the test reagents except the liquid being analyzed are usually incorporated into the element itself. The indicator composition may be dispersed or dissolved in a suitable organic solvent within the element. When the liquid to be analyzed is contacted with the analytical element, some of the liquid is imbibed into the element. Any predetermined analyte present in the imbibed liquid then reacts with oxygen in the presence of the enzyme incorporated in the element to produce a peroxide. The peroxide is formed within the element itself and is situated in close proximity to, or interspersed with, the organic solvent containing indicator composition and a substance having peroxidative activity. It is desirable at this point that all of the peroxide formed in the element should act in the presence of the substance having peroxidative activity to oxidize some of the proximately located indicator composition. This oxidation of indicator composition produces a detectable product whose relative concentration is then determined by measuring its optical density spectrophotometrically or otherwise to indicate the concentration of predetermined analyte in the liquid undergoing analysis. It is apparent that any gentisic acid which is situated in similar proximity to the peroxide as is the indicator composition may itself be oxidized by the peroxide. Any peroxide undergoing such reaction is thus made unavailable for oxidation of indicator composition. What is not so apparent is the reason why significant amounts of gentisic acid come to be as well situated for this reaction as is the indicator composition which is dissolved in the organic solvent. Since gentisic acid is originally dispersed throughout the body of the liquid being analyzed, while indicator composition is in organic solvent within the element itself where the peroxide is first formed, one would expect that most peroxide formed would react with indicator composition before it had a chance to come into contact with significant amounts of gentisic acid. However, this is not the case, and it is accordingly hypothesized that the organic solvents previously chosen to facilitate dispersion of the indicator composition within the element, e.g., a solvent such as N,N-diethyl lauramide, which is used in the prior art (see, for example, U.S. Ser. No. 712,972 referred to above), act also to preferentially partition gentisic acid into the organic solvent from the aqueous liquid. This means that a much higher concentration of gentisic acid may be found in the organic solvent than in the aqueous liquid and thus is just as well situated to react with any peroxide being formed as is the indicator composition itself.\nAccordingly, it would be desirable to provide an analytical element using the peroxide-linked assay mechanism wherein gentisic acid is not preferentially partitioned into the organic solvent containing indicator composition and is thus not as well situated for reaction with peroxide as is the indicator composition and, therefore, does not interfere with the formation of detectable product to such a significant extent as it does in the analytical elements of the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to a frequency divider circuit and a digital phase locked loop (PLL) circuit including the same.\n2. Description of the Related Art\nFIG. 1 is a block diagram of a general programmable digital PLL circuit.\nAs shown in FIG. 1, a digital PLL circuit 6 comprises, for example, a phase comparator 2, a digital counter 8, a frequency multiplier 4, and a frequency divider 5.\nThe phase comparator 2 compares a phase of a reference clock signal of a frequency f.sub.ref with that of an oscillation output f5 from the frequency divider 5 and outputs an up/down signal to a digital counter 8 in accordance with the result of the comparison. For example, when the frequency of the oscillation output f5 is lower than the reference clock signal, it outputs an up signal to the digital counter 8, while in the opposite case, it outputs a down signal to the digital counter 8.\nThe digital counter 8 counts up or counts down the count value from the least significant bit toward the most significant bit based on the up/down signal from the phase comparator 2 and outputs an n-bit count value to the frequency multiplier 4.\nThe frequency multiplier 4 has the same function as a voltage controlled oscillator (VCO), determines an oscillation frequency in accordance with an input count value S3, and finally outputs the target clock S4 of the frequency f.sub.0.\nThe frequency divider 5 outputs an oscillation output f5 obtained by dividing the output signal S4 from the frequency multiplier 4 to the phase comparator 2.\nThe digital PLL circuit 6 shown in FIG. 1 requires an operation time of as much as 2.sup.n /f.sub.ref to reach a locked status shown in FIG. 2 when the digital counter 8 is an n-bit counter.\nIn the digital PLL counter 6, the digital counter 8 is provided with a 32/33 frequency divider which selectively performs frequency division by 32 or 33 and uses this 32/33 frequency divider to count up or count down.\nFIG. 3 is a circuit diagram of a frequency divider 1 of the related art which is provided in the digital counter 8 in FIG. 1.\nFIGS. 7A to 7N and FIGS. 8A to 8N are timing charts of input signals S0, S7, S9, S11, and S14, and frequency division ratio determining signals S21, S14, S17, and S19.\nFIGS. 7A to 7N are timing charts in the case where a 4/5 selecting signal S24 shown in FIG. 3 is a high level (when 4 is selected as a frequency division ratio in the circuit module 3). FIGS. 8A to 8N are timing charts in the case where the 4/5 selecting signal S24 shown in FIG. 3 is a low level (when 5 is selected as a frequency division ratio in the circuit module 3).\nThe frequency divider 1 divides the frequency of an input signal S0 by 32 or 33 in accordance with the 4/5 selecting signal S24.\nAs shown in FIG. 3, the frequency divider 1 comprises the circuit modules 3 and 5.\nThe circuit module 3 comprises D-type flip-flops (D-FFs) 7, 9, and 11, an AND circuit 13, and an OR circuit 14.\nThe D-FFs 7, 9, and 11 are driven using the input signal S0 as a reference clock.\nThe circuit module 3 divides the input signal S0 by 4 or 5 based on the frequency division ratio determining signal S21, shown in FIG. 7J and FIG. 8J, input from the circuit module 5 and outputs the divided signal S7 from a Q.sup.-- terminal of the D-FF 7 to the circuit module 5. Specifically, the circuit module 3 produces the signal S7 shown in FIG. 8B obtained by dividing the input signal S0 by 5 when the frequency division ratio determining signal S21 is a high level, and produces the signal S7 shown in FIG. 7B obtained by dividing the input signal S0 by 4 when the frequency division ratio determining signal S21 is a low level.\nThe circuit module 5 comprises D-FFs 15, 17, and 19, a 4-input NOR circuit 21, and a buffer 23.\nIn the circuit module 5, a CLK terminal of the D-FF 15 is connected to a Q.sup.-- terminal of the D-FF 7 in the circuit module 3, a Q terminal of the D-FF 15 is connected to a CLK terminal of the D-FF 17, and a Q terminal of the D-FF 17 is connected to a CLK terminal of the D-FF 19. Also, in the D-FFs 15, 17, and 19, the D terminals and Q.sup.-- terminals are connected.\nHere, the D-FFs 15, 17, and 19 are connected in series and each D-FF can divide a signal into two. Accordingly, a signal S19 shown in FIGS. 7N and 8N obtained by dividing the signal S7 by 8 (=2.sup.3) is output at the Q terminal of the D-FF 19.\nThe signal S19 is output as an output signal S1 via the buffer 23.\nA signal S15 shown in FIGS. 7L and 8L obtained by dividing the signal S7 by 2 (.dbd.2.sup.1) is output from the Q terminal of the D-FF 15, and a signal S17 shown in FIGS. 7M and 8M obtained by dividing the signal S7 by 4(=22) is output from the Q terminal of the D-FF 17.\nThe NOR circuit 21 receives as input four signals, that is, the signals S15, S17, and S19 from the Q terminals of the D-FFs 15, 17, and 19 and the 4/5 selecting signal S24, and outputs the result of the NOR operation to the AND circuit 13 in the circuit module 13 as a frequency division ratio determining signal S21. Here, the frequency division ratio determining signal S21 becomes a high level, as shown in FIGS. 7J and 8J, when all of the signals S15, S17, and S19 and the 4/5 selecting signal S24 are a low level, while becomes a low level in other cases.\nIn the case of dividing a signal by 32 in the frequency divider 1, the 4/5 selecting signal S24 is held at a high level and the signal S7 obtained by dividing the input signal S0 by 4 is divided by 8 in the circuit module 5. As a result, an output signal S1 obtained by dividing the input signal S0 by 32 is produced.\nOn the other hand, when the frequency divider 1 divides a signal by 33, it makes the circuit module 3 act as a 1/4 frequency divider for seven cycles out of 8 cycles of the signal S7 and act as a 1/5 frequency divider for one cycle out of eight cycles. Due to this, the operation becomes (4.times.7/8+5.times.1/8).times.8, so the frequency divider 1 produces the output signal S1 obtained by dividing the input signal S0 by 33.\nThe problem in the related art was, however, that the PLL circuits used in the cellular phone and other communications fields mainly use frequency dividers containing bipolar, not MOS logic, since the local frequencies have high frequency bandwidths of 1 GHz or more.\nAlso, the power source voltage of PLL circuits used in such communications fields is 3V in most cases, and a basic type of a D-FF has the configuration shown in FIG. 4.\nNamely, the D-FF comprises differential amplifier circuits 200 and 201, emitter-coupled logic (ECL) circuits 202 and 203, and latch circuits 204 and 205.\nThe differential amplifier circuit 200 comprises emitter-coupled npn-type transistors Q1 and Q2 and a constant current source I0 provided at the coupling point. The differential amplifier circuit 201 comprises emitter-coupled npn-type transistors Q3 and Q4 and a constant current source I1 provided at the coupling point.\nThe ECL circuit 202 comprises emitter-coupled npn-type transistors Q5 and Q6. The ECL circuit 203 comprises emitter-coupled npn-type transistors Q9 and Q10.\nThe latch circuit 204 comprises collector-, base-, and emitter-coupled npn-type transistors Q7 and Q8. The latch circuit 205 comprises collector-, base-, and emitter-coupled npn-type transistors Q11 and Q12.\nIn this circuit configuration, the output amplitude of the D-FF can only be about 0.3V or less. It is necessary to reduce the load resistance to improve the through rate.\nHowever, recent cellular phones are expected to provide longer call times, therefore if the load resistance is made small as mentioned above, this will result in an increase of the current consumption and the power consumption.\nAlso, when the through rate is poor, the jitter increases in the output of a bipolar ECL circuit and noise increases in the VCO output signal of the PLL circuit. As a result, the bit error rate of the digital communications signal becomes Inferior.\nFor example, in the D-FF in FIG. 4, when the waveforms of an E input signal and an F input signal produced by an input signal from the D terminal are as shown in FIG. 5A, jitter Ax shown in FIG. 5B is generated in the output signals G and H.\nNote that in the frequency divider 1 shown in FIG. 3, the D-FFs 15, 17, and 19 are serially connected in an asynchronous mode in the circuit module 5.\nAccordingly, the jitter occurring at the D-FF 15 is transmitted to the D-FFs 17 and 19, and jitter .DELTA.Y, which is three times the jitter AX, occurs in the output signals G and H output from the final stage D-FF 19 as shown in FIG. 5C.\nConsequently, in the frequency divider 1 shown in FIG. 3, the jitter becomes large in the finally obtained output signal S1. If the frequency divider 1 is used in a PLL circuit, the phase noise of the VCO output signal of the PLL circuit will increase and the bit error rate of the digital communications signal will end up becoming inferior."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrical connector for being mounted to a circuit board, and more particularly to an electrical connector with an improved spacer for heat dissipation.\n2. Description of Related Art\nWith rapid development of electronic technologies, electrical connectors have been widely used in electronic devices for exchanging information and data with external devices. A conventional QSFP connector usually includes an insulative housing, a plurality of contacts received in the insulative housing, a spacer for organizing the contacts and a metallic shielding cage enclosing the insulative housing. Each contact includes a soldering portion extending beyond the insulative housing for being soldered to a circuit board.\nHowever, since the spacer and the contacts are wholly embedded, the air permeability of spacer of the conventional QSFP connector is poor. As a result, heat generated by the contacts cannot be easily dissipated to the air, thereby decreasing the working life of the QSFP connector.\nHence, an electrical connector with an improved spacer for robust heat dissipation is desired."}
{"text": "This invention relates to the selective conversion of aliphatic and aromatic aminonitriles such as .epsilon.-aminocapronitrile into the corresponding lactams such as .epsilon.-caprolactam by employing a silica catalyst particularly a silica catalyst in the form of spherical beads having high BET surface area, narrow pore and grain size distributions.\nN-substituted amides, especially 5, 6 and 7 membered lactams, are important raw materials for nylon 4, 5 and 6.\nU.S. Pat. No. 2,357,484 (E. L. Martin) discloses a vapor phase process for preparing compounds containing the N-substituted amide group, for example, .epsilon.-caprolactam, by passing a vaporized mixture of water and an aliphatic amino-hydrogen-containing aminonitrile, or a vaporized mixture of water and a nitrile and amino-hydrogen-containing amine over a dehydration catalyst such as activated alumina, silica gel titanium oxide or borophosphoric acid. U.S. Pat. No. 2,357,484 also discloses that diamides are produced by passing a vaporized mixture of water and dinitriles and monoamines over the dehydration catalyst."}
{"text": "Hepatitis B is an infectious disease of the liver caused by the Hepatitis B virus (HBV). The illness can be acute causing liver inflammation, vomiting, jaundice and in some rare instances of severe fulminant disease, death. The majority of infections result in chronic illness that can be either asymptomatic or resulting in chronic liver inflammation leading to cirrhosis of the liver and an increased incidence in the development of hepatocellular carcinoma (HCC).\nHBV infection is a global problem with approximately >350 million people world wide chronically infected and 600,000 die each year from HBV-related liver disease or HCC. The disease has caused epidemics in Asia and Africa and is endemic in China. Transmission of HBV is via infectious blood or body fluids. There are currently vaccines for the prevention of HBV infection. However, the vaccine is prophylatic and cannot be used to treat already infected patients.\nThere are several approved chemotherapeutic treatments for chronic hepatitis B (lamivudine, adefovir, entacavir and telbivudine) that prevent replication of HBV by blocking the action of the HBV polymerase. None of the treatments result in complete clearance of the virus and these treatment result in drug-resistant HBV variants developing after prolonged treatment. Hepatitis B can also be treated with pegylated Interferon-alpha2a but this treatment is associated with severe side effects and is not effective in all patients.\nHBV is a member of the Hepadanvirus family and is divided into 4 major stereotypes (adr, adw, ayr, ayw) based on antigenic epitopes. The virus is also classed into eight genotypes (A-H based on genomic sequence. Genotype A is common in the Americas, Africa, India and Western Europe. Genotype B and C are found in Asia and the US. Genotype D is most common in Southern Europe and India. Genotype E is found in Western and Southern Africa. Genotype F and H are commonly found in Central and Southern America. Genotype G is commonly found in Europe and the U.S.\nHBV genome consists of a circular strand of DNA that is partially double stranded. The genome is ˜3.3 Kb in length. It encodes 4 known genes (C, X, P and S). HBV is one of the few DNA viruses that utilize reverse transcriptase in the replication process. HBV replication involves multiples stages including entry, uncoating and transport of the virus genome to the nucleus. Initially, replication of the HBV genome involves the generation of an RNA intermediate that is then reverse transcribed to produce the DNA viral genome.\nSince current therapies are limited due to their ineffectiveness, serious side effects or due to the generation of drug resistant variants, there is a clinical need for the development of new therapies to treat HBV infection.\nAlteration of viral gene expression, specifically HBV gene expression, through RNA interference (hereinafter “RNAi”) is one approach for meeting this need. RNAi is induced by short single-stranded RNA (“ssRNA”) or double-stranded RNA (“dsRNA”) molecules. The short dsRNA molecules, called “short interfering nucleic acids (“siNA”)” or “short interfering RNA” or “siRNA” or “RNAi inhibitors” silence the expression of messenger RNAs (“mRNAs”) that share sequence homology to the siNA. This can occur via cleavage of the mRNA mediated by an endonuclease complex containing a siNA, commonly referred to as an RNA-induced silencing complex (RISC). Cleavage of the target RNA typically takes place in the middle of the region complementary to the guide sequence of the siNA duplex (Elbashir et al., 2001, Genes Dev., 15:188). In addition, RNA interference can also involve small RNA (e.g., micro-RNA or miRNA) mediated gene silencing, presumably through cellular mechanisms that either inhibit translation or that regulate chromatin structure and thereby prevent transcription of target gene sequences (see for example Allshire, 2002, Science, 297:1818-1819; Volpe et al., 2002, Science, 297:1833-1837; Jenuwein, 2002, Science, 297:2215-2218; and Hall et al., 2002, Science, 297:2232-2237).\nSeveral studies have attempted to use RNAi for the treatment of HBV and this approach has been comprehensively reviewed in RNAI for treating Hepatitis B Viral Infection, Chen, Y. et al., Pharmaceutical Research, 2008 Vol. 25, No. 1, pgs 72-86. Yet, as noted in the above reference transfection of a single siRNA often fails to provide adequate gene silencing. Id. Thus, despite significant advances in the field of RNAi, there remains a need for agents that can effectively inhibit HBV gene expression and that can treat disease associated with HBV expression such as liver disease and cancer."}
{"text": "The present invention relates to novel hydantoin derivatives, processes for producing hydantoin derivatives, pharmaceutical compositions containing at least one of said hydantoin derivatives as aldose reductase inhibitors and novel intermediate compounds in the synthesis of said hydantoin derivatives.\nThe present invention further relates to pharmaceutical compositions containing at least one of hydantoin derivatives as hypoglycemic agents.\nCataract, peripheral neuropathy, retinopathy and nephropathy associated with diabetes mellitus result from abnormal accumulation of polyol metabolites converted from sugars by aldose reductase. For example, sugar cataract results from damage of lens provoked by change in osmotic pressure induced by abnormal accumulation of polyol metabolites converted from glucose or galactose by aldose reductase in lens. [see J. H. Kinoshita et al., Biochim. Biophys. Acta, 158, 472 (1968) and cited references in the report]. And some reports were submitted about undesirable effect of abnormal accumulation of polyol metabolites in lens, peripheral nerve cord and kidney of the diabetic animals [see A. Pirie et al. Exp. Eye Res., 3, 124 (1964) ; L. T. Chylack Jr. et al., Invest. Ophthal., 8, 401 (1969) J. D. Ward et al., Diabetologia, 6, 531 (1970)]. Consequently, it is important to inhibit aldose reductase as strongly as possible for treating and/or preventing diabetic complications mentioned above. Although several compounds have been offered as aldose reductase inhibitors, none of them is fully sufficient in inhibitory activity against the enzyme. Therefore, it has been desired to develop new compounds having a stronger inhibitory activity against aldose reductase.\nIn spite of the early discovery of insulin and its subsequent wide-spread use in the treatment of diabetes mellitus, and the later discovery and use of sulfonylureas (e.g. chlorpropamide, tolbutamide) and biguanides (e.g. phenformin) as oral hypoglycemic agents, the treatment of diabetes mellitus remains less than satisfactory. Insulin can only be administered intravenously due to its chemical nature, and therefore, is troublesome and inconvenient to use. Oral hypoglycemic agents tend to promote side effects such as excessive hypoglycemia or lactic acidosis. A continuing need for potent hypoglycemic agents, which may be less toxic, is clearly evident.\nRecently, developments of the aldose reductase (AR) inhibitors as agents for diabetic complications are in progress. AR inhibitors will not be agents for diabetes mellitus, but symptomatic agents for diabetic complications, so they are expected to be little effective against diabetes mellitus itself.\nDiabetes mellitus should be treated by hypoglycemic agents, and preferably, by the hypoglycemic agents with AR inhibiting activity, so such agents having both hypoglycemic and AR inhibiting activities have been desired.\nFurther, diabetes mellitus is usually accompanied by cardiovascular disease due to atherosclerosis, so the hypoglycemic agents with hypolipidemic activities have also been desired."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of driving a gate line, a gate drive circuit and a display apparatus having the gate drive circuit. More particularly, the present invention relates to a method of driving a gate line in which reliability of an operation thereof is substantially enhanced, a gate drive circuit performing the method and a display apparatus having the gate drive circuit.\n2. Description of the Related Art\nGenerally, a liquid crystal display (“LCD”) apparatus includes an LCD panel which displays images by controlling a light-transmitting ratio of liquid crystal molecules provided with light from a backlight assembly disposed below the LCD panel.\nThe LCD apparatus typically includes a display panel, a gate driving part and a data driving part. The display panel includes gate lines, data lines and pixel parts electrically connected to the gate lines and the data lines. The gate driving part outputs gate signals to the gate lines. The data driving part outputs data signals to the data lines. The gate driving part and the data driving part may be formed in a chip mounted on the display panel.\nIn efforts to decrease a size of the LCD apparatus and to improve manufacturing productivity of the LCD apparatus, the gate driving part has been integrated on a display substrate such as an amorphous silicon gate (“ASG”) type substrate. However, when a gate drive circuit integrated on the display substrate is driven at a high temperature, noise from an abnormal gate on signal is generated during a gate off signal interval is generated.\nIn addition, properties of an amorphous silicon thin-film transistor (“a-Si TFT”) of the gate drive circuit vary with changes in temperature and over time.\nAs a result, a gate bias stress applied to a TFT in the gate drive circuit generates a threshold voltage shift when the gate drive circuit is driven. The threshold voltage shift substantially reduces a current driving ability of a hold transistor which is included in the gate drive circuit to maintain an off level of a gate line and a gate terminal of a pull-up transistor which outputs a gate signal."}
{"text": "For the present purposes, \"welding\" can fairly be said to embrace the fusing together of, in particular, two or more metal parts. Welded joints often exhibit greater rigidity than bolted or riveted constructions, and good quality welds are nonporous and leak-proof. In point of fact, good quality welds approximate the level of strength of the parent materials. Weld strengths are only slightly reduced below those levels as a consequence of unfavourable heat stresses that are set up in the immediate vicinity of the weld during the welding process. The potential advantages of welding in manufacturing processes include reduced capital investment requirements, greater flexibility of design, quicker change over to new or alternate designs, reduced machining and cleaning of parts, and improved strength to weight ratios in the assemble product.\nThe realization of any one or more of these associated advantages is, of course, contingent on the quality of the weld in question. There are a number of factors which impact on the character and quality of the weld, including the selection of any given welding technology, the competency of the operator, and of particular importance in the present context, the condition of the welding equipment.\nA commonly used welding technique involves resistance welding. At least two welding electrodes are arranged in mutually opposed relation along a common axis, along which they are relatively movable. The two or more work pieces that are sought to be welded ar interposed between the two electrodes while they are arranged in axially spaced relation from one another. When the workpieces are properly mutually aligned there between, the electrodes are moved towards one another, and embrace the workpieces in forcibly clamped relation, whereupon an electric current is passed between the electrodes and the heat generated by the electrical resistance of the interposed workpieces results in localized melting of the workpieces proximal to the contacting electrode surfaces. The melted materials from the two or more clamped workpieces fuses together and the intermingled materials harden into a unified piece once the electrical current is discontinued. This welding technology utilizes large, physically robust electrodes in order to provide the prerequisite clamping strength. Typical electrodes used in industrial applications may be one half inch in diameter, or more.\nThe quality of the weld is to some degree, contingent on the condition of the mating surfaces of the electrodes. It was with this in mind that a variety of devices were produced, which were intended to recondition the electrodes. Such a treatment is necessary since the surfaces of the electrodes degrade quite quickly over the course of normal use. Examples of surface reconditioning apparatus for use in treating resistance welding electrodes are disclosed in the following patents: U.S. Pat. Nos. 4,682,487; 4,856,949; 4,916,931; and 4,921,377.\nARC welding is another well known welding technique. This differs fundamentally from resistance welding in that ARC welding electrodes are deliberately consumed during the welding process, and thereby come to form an integral component of the welded product. Accordingly, the problem of electrode reconditioning that arises in association with resistance welding, does not arise in ARC welding practices.\nMIG (metal-inert-gas arc) welding is such an ARC welding process. More particularly, it is a process in which the electrode, in the form of a relatively fine wire, is continuously fed from a large spool driven by a variable speed welding drive. The speed at which the wire electrode is delivered to the weld is controlled in order to optimize arc length and burnoff rate during the welding process.\nThe electrical arc is enveloped in a moving gas flow, usually argon or other inert gas, or mixtures thereof. In an especially preferred form, MIG welding utilizes carbon dioxide as a shielding gas.\nIn MIG welding generally, both the wire electrode and the gas are channelled through a so-called \"torch\", which includes a central, electrically charged \"tip\". The tip directs the wire electrode toward the weld site, and a concentrically arranged metal gas shield that is electrically insulated from the tip, acts as a hood to direct and maintain a coaxial flow of the inert gas in surrounding relation about the wire. The quality of the weld is contingent on both consistent and continuous gas flow and arc patterning. Anything which interferes with the gas flow or redirects or otherwise militates against the desired electrical arc pattern, will diminish the quality of the weld.\nMIG welding, when properly executed, permits high welding speeds, and allows for less operator training than is required in the case of other welding techniques. In applications where one or the other or both of these benefits are sought, the weld quality is especially sensitive to those variations which are attributable to adverse gas flow or arc patterning influences.\nGas flow in MIG welding can be adversely effected as a consequence of molten metal deposition. This arises as a result of backsplash splatter on the respective mutually opposed surfaces of the tip and the hood, within the interior of the torch.\nSimilarly, (since the dielectric strength of the gas flow is otherwise a constant), the accumulation of such backsplash splatter decreases the physical and hence \"electrical\" distance between the charged tip and the electrically insulated hood. If the distance decreases sufficiently, the voltage differential will exceed the dielectric strength of the intervening gas flow, and the arc will jump between the tip and the hood. This results in a diminished amount of electrical energy being delivered to the weld site and a concomitant compromise in weld quality.\nIn view of the foregoing, it is important that MIG welding torches be cleaned regularly, in order to avoid these two latter mentioned problems. This realization has led to the development of a number of devices that are intended to perform the necessary operations.\nBy way of example, one such device, which is intended for use in robotic MIG welding operations, there is provided a heavy gauge wire that is clamped at one end thereof, in upstanding relation, with its uppermost free end available to be received internally of the torch, between the hood (or gas shield) and the tip. The robotic arm is preprogrammed to essay the torch along a predetermined circular path during the cleaning cycle, so that the upstanding wire dislodges splatter material from the two opposed surfaces of the hood and the tip. This approach to the problem can result in the tip being bent out of concentric alignment within the hood, which will in turn result in the very problems that the cleaning cycle is intended to help avert.\nTwo other such devices each employ a two-part clamping chuck, having a fixed jaw and a movable jaw. Such arrangements do not compensate for differences or vacations in torch nozzle sizes, or off-centred torch insertion, and can result in significant torch nozzle damage.\nThese clamps are intended to secure the torch in a rigidly-held and centred position, relative to an axially aligned rotating platform. This platform supports one end of each of a number of rotatable, longitudinally extending, substantially elongated blades which are aligned in such a way as to extend in free standing relation within the space between the hood and the tip, and upon rotation to dislodge the splatter from the two surfaces. In both such devices the design of the clamping chuck with its stationary jaw and the use of the substantially elongated blades, can still result in the tip being bent out of alignment relative to the hood, with the seriously adverse consequences already alluded to herein before.\nAs a consequence of the forgoing, there remains a need in the art for devices that are adapted to minimize the risk of misaligning the tip, while at the same time effectively removing the splatter from within the MIG torch."}
{"text": "Some homes today are equipped with smart home networks to provide automated control of devices, appliances and systems, such as heating, ventilation, and air conditioning (“HVAC”) system, lighting systems, alarm systems, home theater and entertainment systems. Smart home networks may include control panels that a person may use to input settings, preferences, and scheduling information that the smart home network uses to provide automated control the various devices, appliances and systems in the home. For example, a person may input a desired temperature and a schedule indicating when the person is away from home. The home automation system uses this information to control the HVAC system to heat or cool the home to the desired temperature when the person is home, and to conserve energy by turning off power-consuming components of the HVAC system when the person is away from the home. Also, for example, a person may input a preferred nighttime lighting scheme for watching television. In response, when the person turns on the television at nighttime, the home automation system automatically adjusts the lighting in the room to the preferred scheme."}
{"text": "Ad-hoc networks include multiple devices or nodes that exchange wireless signals. During operation, nodes may enter and leave the proximity of other nodes. Thus, the composition of an ad-hoc network may change over time. Moreover, the mobility of nodes may cause changes in various network characteristics, such as topology. Despite a lack of centralized authority or existing infrastructure, ad-hoc networks are typically capable of rearranging themselves in response to such events.\nRecently, ad hoc networking techniques have been considered an attractive technology for implementing mesh networks, which provide a multipoint-to-multipoint network topology. In such networks, communication between two devices may occur across one or more intermediate or relaying nodes. Such communications are referred to as multihop communications.\nThe application of ad-hoc communications techniques to multihop networking is viewed as a way to provide new applications for mobile device users. In addition, this application has the potential to provide new opportunities for the communications industry in the areas of terminal manufacturing, software engineering, and the deployment of network infrastructure to interconnect ad-hoc networks. Moreover, this application of ad-hoc communications techniques to multihop networking provides for various consumer uses. Examples of such uses include applications related to teenager and other group networking, Internet access, authentication applications, and home networking.\nBluetooth and wireless local area networks (WLAN) are examples of wireless ad-hoc networking technologies. Bluetooth provides a short-range radio network, originally intended as a cable replacement. It can be used to create ad hoc networks of up to eight devices, where one device is referred to as a master device. The other devices are referred to as slave devices. The slave devices can communicate with the master device and with each other via the master device. The devices operate in the 2.4 GHz radio band reserved for general use by Industrial, Scientific, and Medical (ISM) applications. Bluetooth devices are designed to find other Bluetooth devices within their communications range and to discover what services they offer.\nWLANs are local area networks that employ high-frequency radio waves rather than wires to exchange information between devices. IEEE 802.11 refers to a family of WLAN standards developed by the IEEE. In general, WLANs in the IEEE 802.11 family provide for 1 or 2 Mbps transmission in the 2.4 GHz band using either frequency hopping spread spectrum (FHSS) or direct sequence spread spectrum (DSSS) transmission techniques.\nWithin the IEEE 802.11 family are the IEEE 802.11b and IEEE 802.11g standards. IEEE 802.11b (also referred to as 802.11 High Rate or Wi-Fi) is an extension to IEEE 802.11 and provides for data rates of up to 11 Mbps in the 2.4 GHz band. This provides for wireless functionality that is comparable to Ethernet. IEEE 802.11b employs DSSS transmission techniques. IEEE 802.11g provides for data rates of up to 54 Mbps in the 2.4 GHz band. For transmitting data at rates above 20 Mbps, IEEE 802.11g employs Orthogonal Frequency Division Multiplexing (OFDM) transmission techniques. However, for transmitting information at rates below 20 Mbps, IEEE 802.11g employs DSSS transmission techniques. The DSSS transmission techniques of IEEE 802.11b and IEEE 802.11g involve signals that are contained within a 23 MHz wide channel. Several of these 23 MHz channels are within the ISM band.\nOther technologies are also applicable for the exchange of information at higher data rates. Ultra wideband (UWB) is an example of such a higher data rate technology. Since gaining approval by the Federal Communications Commission (FCC) in 2002, UWB techniques have become an attractive solution for short-range wireless communications. Current FCC regulations permit UWB transmissions for communications purposes in the frequency band between 3.1 and 10.6 GHz. However, for such transmissions, the spectral density has to be under −41.3 dBm/MHz and the utilized bandwidth has to be higher than 500 MHz.\nThere are many UWB transmission techniques that can fulfill these requirements. A common and practical UWB technique is called impulse radio (IR). In IR, data is transmitted by employing short baseband pulses that are separated in time by gaps. Thus, IR does not use a carrier signal. These gaps make IR much more immune to multipath propagation problems than conventional continuous wave radios. RF gating is a particular type of IR in which the impulse is a gated RF pulse. This gated pulse is a sine wave masked in the time domain with a certain pulse shape.\nTo participate in an ad-hoc multihop network, a device needs to provide several features. Examples of such features include interference avoidance, link management, and routing. Moreover, certain wireless communication technologies are better suited for the exchange of control information, while other wireless communication technologies may be better suited for the transfer of user data. For instance, Bluetooth on its own is not well suited for many forms of user data. However, higher data rate technologies (e.g., WLAN and UWB) are often not efficient for the transfer of network control information.\nTherefore, techniques are needed for the effective use of ad hoc techniques. in multihop networks."}
{"text": "1. Field of the Invention\nThe present invention relates to a garnet-type ion conducting oxide, a complex, a lithium secondary battery, a manufacturing method of a garnet-type ion conducting oxide and a manufacturing method of a complex.\n2. Description of the Related Art\nGarget-type oxides such as Li7La3Zr2O12 and Li7ALa3Nb2O12 (A=Ca, Sr or Ba) synthesized by the solid-phase reaction method have been proposed conventionally as a solid electrolyte configured to conduct lithium ion (Non-Patent Literatures 1 to 3). It has been reported that this solid electrolyte has the conductivity of 1.9 to 2.3×10−4 Scm−1 (25° C.) and activation energy of 0.34 eV. The inventors have studied a solid electrolyte of Li7La3Zr2O12-based garnet-type ion conducting oxide among garnet-type oxides having excellent chemical stability and a wide potential window. For example, it has been proposed that the Zr sites in this solid electrolyte should be substituted with an element such as Nb, in order to enhance the conductivity (see, for example, Patent Literature 1). This solid electrolyte has high conductivity but needs treatment at high temperature such as 1200° C. It has been proposed, on the other hand, that La sites should be additionally substituted with an alkaline-earth metal, in order to minimize reduction of the electric conductivity and reduce the firing energy (see, for example, Patent Literature 2).\nA solid electrolyte including Li, La, Zr, O and Al has been proposed as Li7La3Zr2O12-based solid electrolyte (see, for example, Patent Literature 3). According to the disclosure of this prior art, addition of Al to Li7La3Zr2O12-based solid electrolyte provides the solid electrolyte with the density and the conductivity required for the solid electrolyte material."}
{"text": "This invention relates to a fluid level indicator for a small watercraft, and more particularly to a fluid level indicator which can provide a useful indication of a particular fluid level in a small watercraft even when the watercraft is experiencing severe rocking and up and down movements during operation and which is able to detect the fluid level at a number of positions.\nOne type of particularly popular small watercraft is of the jet propelled type and is designed to be operated by a single rider who is seated on the seat in straddle-like fashion. This type of small watercraft is highly maneuverable and is very sporting in nature. A control bridge is normally located forwardly of the seat and carries handlebars which are positioned where the rider may conveniently grasp them to steer the watercraft. An instrument panel may also be positioned on the control bridge in view of the rider and may include an indicator for a fluid such as fuel or oil.\nOne type of known device for detecting the fluid level in a fuel or oil tank of a boat includes a pivotal arm having a float at one end thereof. Fluctuation of the fluid level causes vertical displacement of the float so that the fluid level may be detected in reference to the angular orientation of the pivotal arm. While this type of device is generally satisfactory for large watercraft, it is not well suited for smaller watercraft which are more strongly influenced by wind and water forces. When this type of device is applied to a small watercraft, the float tends to vibrate or move when the fluid is agitated as a result of these wind and/or water forces which cause the small watercraft to rock or move up and down.\nAnother type of device used for detecting the fluid level in a fuel or oil tank of a boat includes a vertical guide member disposed in the tank, a single float vertically movable and slidably supported on the guide, and reed switches positioned within the guide for detecting the position of the float along the guide. FIG. 1 illustrates such a device. In this figure, the vertical guide member is designated by the numeral 61 within which three vertically spaced apart reed switches 62a through 62c are disposed. An annular float 63 having a permanent magnet 64 is slidably received by the guide 61. The switches 62a through 62c are turned on and off on based on the position of the float 63 in the sections A through E as follows:\n______________________________________ POSITION OF SWITCH THE FLOAT 62a 62b 62c ______________________________________ Section E off off on Section D off off off Section C off on off Section B off off off Section A on off off ______________________________________\nAs can be seen from the above, this type of arrangement does not distinguish between the float being in section B and in section D. This problem still occurs even if the number of switches are increased to four or more. Decreasing the number of switches to one or two with this type of arrangement will decrease the number of sections so that the fluid level can be detected at only a limited number of positions.\nIt is therefore a principal object of this invention to provide a fluid level indicator which can be used to effectively detect and indicate a fluid level in a small watercraft even when the watercraft is experiencing severe rocking or up and down movements during operation.\nIt is a further object of this invention to provide a fluid level indicator for a small watercraft which is able to detect and distinguish the fluid level at a number of different positions."}
{"text": "Conventionally, an optical waveguide structure in which an optical waveguide configured to be coupled to an optical fiber is layered on a substrate formed with a groove for positioning the optical fiber, and an optical-waveguide-type optical module including the optical waveguide structure have been known (see the Patent Publication 1 below).\nFurther, conventionally, an optical waveguide structure which is a component of an optical-waveguide-type optical module, an optical-waveguide-type optical module, and an optical fiber array configured to be coupled to an optical waveguide so as to form an optical module have been known. Particularly, the optical waveguide structure, the optical-waveguide-type optical module, and the optical fiber array, each of which having a groove for supporting an optical fiber, have been known (see the Patent Publication 1 below).\nReferring to FIGS. 17-19, a first example of a conventional optical waveguide structure will be explained. FIG. 17 is a top plan view showing an optical-waveguide-type optical module including a conventional optical waveguide structure. FIG. 18 is a fragmentary enlarged cross-sectional view taken along the line XVIII-XVIII in FIG. 17, and FIG. 19 is a fragmentary enlarged cross-sectional view taken along the line XIX-XIX in FIG. 17.\nAs shown in FIGS. 17-19, an optical-waveguide-type optical module 200 having a single upstream optical fiber 202 extending longitudinally, eight downstream optical fibers 204 spaced longitudinally from the upstream optical fiber 202 and arranged laterally relative to each other, and an optical waveguide structure 206 for transmitting light through the upstream optical fiber 202 to the downstream optical fibers 204. The upstream optical fiber 202 and the downstream optical fibers 204 include respective cores 202a, 204a extending longitudinally.\nThe optical waveguide structure 206 has a substrate 212 on which a single upstream groove 208 extending longitudinally and eight downstream grooves 210a-210h extending longitudinally and spaced longitudinally from the upstream groove 208 are provided, and an optical waveguide 214 layered on the substrate 212 between the upstream groove 208 and the downstream grooves 210a-210h. The upstream optical fiber 202 is positioned on the upstream groove 208, and the downstream optical fibers 204 are positioned on the downstream grooves 210a-210h. \nThe optical waveguide 214 includes a lower cladding 214a layered on the substrate 212, a core 214b formed on the lower cladding 214a, and an upper cladding 214c layered on the lower cladding 214a and the core 214b. The core 214b of the optical waveguide 214 is formed so that, when the optical fibers 202, 204 are supported and positioned on the upstream groove 208 and the downstream grooves 210a-210h, the core 214b of the optical waveguide 214 is aligned with the cores 202a, 204a of the optical fibers 202, 204 at the same level in an vertical direction.\nFurther, for light transmission between the upstream optical fiber 202 positioned on the upstream groove 208 and each of the downstream optical fibers 204 positioned on the downstream grooves 210a-210h, the core 214b of the optical waveguide 214 has a single upstream port 220 aligned with the upstream groove 208 and eight downstream ports 222 each aligned with the respective downstream grooves 210a-210h. In the illustrated optical waveguide structure 206, the core 214b of the optical waveguide 214 extends from the single upstream port 220, is branched toward a downstream side, and terminates at the eight downstream ports 222. The optical waveguide 214 has an upstream portion 224a adjacent to the upstream port 220, an intermediate portion 224b between the upstream port 220 and the downstream ports 222, and a downstream portion 224c adjacent to one of the downstream ports 222.\nLight transmitted through the single upstream optical fiber 202 is transmitted from the upstream port 220 to the optical waveguide 214 and branched toward the downstream side. Then, the branched lights are transmitted from the eight downstream ports 222 to the eight downstream optical fibers 204. In this case, the optical-waveguide-type optical module 200 serves as an optical splitter. On the contrary, when light is traveled in the opposite direction from the downstream optical fibers 204 to the upstream optical fiber 202, the optical-waveguide-type optical module 200 serves as an optical coupler.\nReferring to FIGS. 20-22, a second example of the conventional optical waveguide structure which is a component of an optical-waveguide-type optical module will be explained. FIG. 20 is a top plan view showing an optical-waveguide-type optical module including the conventional optical waveguide structure. FIG. 21 is a fragmentary enlarged cross-sectional view taken along the line XXI-XXI in FIG. 20, and FIG. 22 is a fragmentary enlarged cross-sectional view taken along the line XXII-XXII in FIG. 20.\nAs shown in FIGS. 20-22, an optical-waveguide-type optical module 300 has a single upstream optical fiber 302 extending longitudinally, eight downstream optical fibers 304 spaced longitudinally from the upstream optical fibers 302 and arranged laterally relative to each other, and an optical waveguide structure 306 for supporting the upstream optical fiber 302 and the downstream optical fibers 304 and transmitting light through the single upstream optical fiber 302 to the eight downstream optical fibers. The optical-waveguide-type optical module 300 further has two fiber-holding lids 308a, 308b which respectively hold the upstream optical fiber 302 and the downstream optical fibers 304 against the optical waveguide structure 306, and an adhesive 310 filled between any two of the optical fibers 302, 304, the optical waveguide structure 306 and the fiber-holding lids 308a, 308b to fix them relative to each other.\nThe upstream optical fiber 302 and the downstream optical fibers 304 have respective cores 302a, 304a extending longitudinally. The optical waveguide structure 306 has a substrate 312 and an optical waveguide 314 layered thereon. The substrate 312 has an upper surface 316 with a lateral width W, and a plurality of grooves 318 for supporting the upstream optical fiber 302 and the downstream optical fibers 304 are formed in the upper surface 316. The optical waveguide 314 includes a core 314a which is formed so that, when the optical fibers 302, 304 are supported and positioned on the grooves 318, the core 314 of the optical waveguide 314 is aligned with the cores 302a, 304a of the optical fibers 302, 304. Each of the fiber-holding lids 302, 304 has the same width as that of the substrate 312, and is provided with (a) contact groove(s) 322 contacting with the optical fiber(s) 302, 304.\nLight transmitted through the single upstream optical fiber 302 is transmitted to the optical waveguide 314, and branched toward a downstream side. Then, the branched lights are transmitted to the eight downstream optical fibers 304. In this case, the optical-waveguide-type optical module 300 serves as an optical splitter. On the contrary, when light is traveled in the opposite direction from the downstream optical fibers 304 to the upstream optical fiber 302, the optical-waveguide-type optical module 300 serves as an optical coupler.\nWhen light is transmitted from the upstream optical fiber 302 to the optical waveguide 314 and transmitted from the optical waveguide 314 to the downstream optical fibers 304, a loss of optical power to be transmitted, called an insertion loss, is caused.\nPatent Publication 1: Japanese Patent Laid-Open Publication No. 11-125731\nIn the above optical-waveguide-type optical module 200 which is the first conventional example, when light is transmitted from the upstream optical fiber 202 to the optical waveguide 214 and transmitted from the optical waveguide 214 to the downstream optical fibers 204, a loss of optical power to be transmitted, called an insertion loss, is caused. The insertion loss is a ratio of a downstream output optical power (Po) relative to an upstream input optical power (Pi) expressed in deci Bell unit, i.e., (10 log10 (Po/Pi)). An amount of the insertion loss of the optical-waveguide-type optical module 200 is preferable as small as possible. In a case of the optical module according to the present invention where there is no gain such as amplifying performance, i.e., Po<Pi, a value of the insertion loss expressed in deci Bell unit based on the above formula is negative. When there is no insertion loss, the value thereof based on the above formula is 0 (zero). Thus, in this specification, a small insertion loss means that a negative value calculated by the above formula is large, namely, an absolute value thereof is close to 0 (zero), as with common practice. Such insertion loss can be applied to a case including both of an insertion loss and a gain. In this case, the insertion loss can be considered after a part of the gain is separated. Therefore, an optical waveguide structure is needed, which structure can reduce an amount of insertion loss of an optical-waveguide-type optical module thereof as compared with the conventional optical waveguide structure.\nAn amount of insertion loss of the optical-waveguide-type optical module 200 shown in FIGS. 17-19, i.e., that of insertion loss between the upstream optical fiber 202 and the optical waveguide 214 or that of insertion loss between the optical waveguide 214 and the downstream optical fibers 204, may be changed depending on changes in an environmental temperature. Particular, since an optical splitter and an optical coupler which are examples of the optical-waveguide-type optical module are used as a part of the Internet with an optical fiber network disposed outside, an environmental temperature thereof may vary, for example, from −40° C. to +85° C. When insertion loss of such an optical-waveguide-type optical module is fluctuated due to variations in the environmental temperature, the optical module may not be able to give its original performance. The fluctuation of the insertion loss of the optical-waveguide-type optical module due to the change in the environmental temperature is preferably as small as possible. Thus, it is desirable to provide an optical waveguide structure capable of reducing fluctuation of an insertion loss of an optical-waveguide-type optical module as compared with that of the conventional optical-waveguide-type optical module.\nIn the above optical-waveguide-type optical module 300 which is the second conventional example, the upstream fiber-holding lid 308a is likely to be fixed to the optical fiber in a state laterally inclined relative to the substrate 312 (see FIG. 22). In this case, on the opposite sides of the upstream optical fiber 302, distances between the fiber-holding lid 308a and the upper surface 316 of the substrate 312 are different from each other as well as thicknesses of the adhesive 310 are different from each other.\nAs mentioned above, an optical splitter and an optical coupler, which are examples of the optical-waveguide-type optical module 300, are used as a part of the Internet with the optical fiber network placed outdoor. An environmental temperature around the optical module 300 may be changed, for example, from −40° C. to +85° C. When the environmental temperature is changed, each of the substrate 312, the fiber-holding lid 308a and the adhesive 310 expands and contracts according to their thermal expansion coefficients which are different from each other. Thus, when thicknesses of the adhesive 310 are different from each other on the opposite sides of the upstream optical fiber 302 and the environmental temperature is changed, the upstream optical fiber is subjected to uneven stress, resulting in deterioration of an insertion loss. Fluctuation of the insertion loss of the optical-waveguide-type optical module 300 when the environmental temperature is changed is preferably as small as possible.\nThus, there is a need to provide an optical waveguide structure which is a component of an optical-waveguide-type optical module and is capable of reducing fluctuation of an insertion loss thereof when an environmental temperature is changed. Further, there is a need to provide an optical-waveguide-type optical module capable of reducing fluctuation of an insertion loss thereof when an environmental temperature is changed.\nFurthermore, there is a need to provide an optical fiber array configured to be coupled to an optical waveguide so as to form an optical module, the optical fiber array being capable of reducing fluctuation of an insertion loss of the optical module when an environmental temperature is changed.\nIt is therefore a first object of the present invention to provide an optical waveguide structure capable of reducing an insertion loss of an optical-waveguide-type optical module including the optical waveguide structure as compared with the conventional optical waveguide structure, and to provide the optical-waveguide-type optical module itself.\nFurther, it is a second object of the present invention to provide an optical waveguide structure capable of reducing fluctuation of an insertion loss of an optical-waveguide-type optical module including the optical waveguide structure when an environmental temperature is changed, and to provide the optical-waveguide-type optical module itself.\nFurther, it is a third object of the present invention to provide an optical fiber array configured to be coupled to an optical waveguide so as to form an optical module, the optical fiber array being capable of reducing fluctuation of an insertion loss of the optical module when an environmental temperature is changed."}
{"text": "1. Field\nThe following description relates to a display apparatus and method for questions and answers, and more particularly, to a display apparatus and method that provide a question and answer service related to non-sentence speech.\n2. Description of the Related Art\nBecause display apparatuses have been capable of recognizing user's voice, display apparatuses have been capable of recognizing user's speech voice to provide a voice recognition result intended by users or perform a control command.\nIn particular, the most universal method related to a voice conversation service that provides the voice recognition result intended by users is a keyword recognition method that performs a conversation between users and display apparatuses through a search engine. Such a keyword recognition method refers to display apparatuses that perform a search related to a core keyword based on speech voice related to the core keyword and provide a search result. However, such a method provides numerous search results related to the speech voice for the core keyword, and thus users experience inconvenience by having to search for desired information from numerous search results.\nAnother method is a sentence recognition method, where display apparatuses analyze speech voice related to a question sentence of a user to determine a user's speech intention, perform a search based on the determined speech intention, and provide an answer result according to the search.\nSuch a sentence recognition method is advantageous compared to the keyword recognition method described above in that the sentence recognition method provides an answer result closer to the user's speech intention. However, when a sentence speech appropriate to a preset pattern is not input, or noise occurs in the sentence speech due to an external environment, such a sentence recognition method does not perform correct voice recognition on the sentence speech of a user."}
{"text": "1. Field of the Invention\nThe present invention relates to the test techniques and testability architectures used in testing integrated circuits. More specifically, the invention pertains to test techniques applied to testing of user-configurable arrays before they are configured by the user.\n2. The Prior Art\nUser-configurable gate arrays consist of logic circuits or blocks that can be connected together by configurable interconnections, such as anti-fuse elements, to implement a desired circuit function. The configurable interconnect consists of interconnect layers such as metallization, and configurable devices which, when programmed, establish electrical connections between the interconnect layers. However, before configuring the circuit to implement a particular function, all the individual modules in the array and all the input/output (I/O) modules and buffers are isolated from one another. This presents a challenging test problem.\nBefore the circuit is configured by the user, all of the active circuits in such integrated circuits such as logic modules, I/O modules, configuring circuits, etc., must be tested and guaranteed to be fully functional and meet all required specifications. In addition, all passive interconnect circuits such as metallization interconnect, anti-fuse elements, feed-thru pass transistors, must also be free of defects and guaranteed. This is necessary so that a customer configuring such a circuit can expect a fully functional, high quality integrated circuit after his application circuit is mapped into the device. It is thus imperative that test architectures and test techniques be developed to solve this problem, namely, how to guarantee full functionality and spec of a one-time programmable user configurable array circuit before the circuit is configured by the user.\nUser configurable arrays or PLDs (programmable logic devices) which use erasable elements to implement their interconnect do not have to contend with this problem since the array can be configured to implement any circuit pattern, be fully tested and later erased to the \"blank\" state for reconfiguration."}
{"text": "1. Field of the Invention\nThis invention relates to adhesive and sealant compositions wherein one or more of the reactive and/or curable or polymerizable constituents of said compositions are encapsulated. It more particularly relates to encapsulated cure systems for, directly or indirectly, initiating and/or effectuating the cure or polymerization of adhesive and sealant compositions, as well as methods of making the same.\n2. Description of Related Art\nThe use of microencapsulated ingredients in the manufacture of pharmaceuticals, pesticides, paints, adhesives, sealants and printing inks is well known. Perhaps the most widely known use of microcapsules has been in the product generally referred to as carbonless paper which incorporate or have applied thereto coatings comprising microencapsulated inks and/or other color forming agents that are released upon the application of pressure to the paper surface. In the manufacturing and service sectors one of the most widely known use of microcapsules has been in the manufacture of adhesives and sealants. Here the use of microencapsulation allows one to form a one-part adhesive or sealant from a typical two- or more-part adhesive and sealant system. Microencapsulation also allows one to pre-apply such adhesives and sealants to substrates at the site of manufacture or conversion rather than at the site of use or installation. The use of microencapsulation in adhesive and sealants is well known and takes many different paths.\nEncapsulated solvent-based adhesive systems are of several different constructs. Roesch et. al. (U.S. Pat. No. 5,922,798) teach solvent based adhesive compositions wherein a solvent, alone or together with a therein dissolved resin, is encapsulated and used to bond two substrates, each of which is dissolved or softened by the solvent. Eichel (U.S. Pat. No. 2,907,682) discloses adhesive tapes wherein the adhesive portion comprises a combination of encapsulated solvent and encapsulated solid adhesives, which are soluble in the solvent. Upon application of the tape to a substrate and the further application of pressure, the capsules rupture allowing the solvent to dissolve or at least tackify the adhesive whereby, upon evaporation or absorption of the solvent, an adhesive bond is formed. Where the solvent is non-evaporating, e.g., a plasticizer, and neither the substrate carrying the adhesive nor the substrate to which it is to be bonded absorb the non-volatile solvent, the result of the combination of the solvent and adhesive material is a pressure sensitive adhesive. Adhesives based on polyvinyl acetate, rubber, nitrile rubber, ethyl cellulose, or other cellulose derivatives such as cellulose acetate are especially suited for the solvent activation/reactivation type application.\nHot melt adhesive systems employing encapsulated solvents are also known. Baetzold et. al. (U.S. Pat. No. 6,084,010) teaches solid, tacky or non-tacky hot melt glue compositions having incorporated therein microcapsules of a solvent capable of softening or further tackifying the hot melt. Rubbing the hot melt composition, typically in the form of a stick, on the substrate to which it is to be applied fractures the capsules, releasing the solvent which, in turn, softens the hot melt and allows it to be deposited on the surface of the substrate.\nAnother type of encapsulated adhesive and sealant is that where the adhesive or sealant material or, in the case of a curable adhesive or sealant material, the components thereof are encapsulated in a single capsule. These capsules are typically applied to a substrate in a binder system that is non-tacky and dry to the touch. In this way, otherwise tacky or liquid flowable adhesives can be pre-applied, but not activated or bond forming until the capsule walls themselves are fractured releasing or exposing the adhesive materials. For example, Eichel (U.S. Pat. No. 2,986,477) teaches the encapsulation of tacky adhesive materials. Wallace (U.S. Pat. No. 4,428,982) teaches the encapsulation of curable anaerobic adhesives wherein the encapsulating material is air permeable so that the curable adhesive remains in a liquid or uncured state in the capsule until use. Schwantes (U.S. Pat. No. 6,592,990) teaches encapsulated adhesives, particularly pressure sensitive adhesives, wherein the adhesive is formed in-situ, after encapsulation of the ingredients therefore.\nA third, and perhaps the most common use of encapsulation in adhesives and sealants, involves curable or reactive adhesive and sealant compositions which rely upon the presence of curatives or curing and/or cross-linking agents and/or other reactants such as activators, catalysts, initiators, accelerators, and the like for effecting polymerization or curing of the composition so as to form the desired adhesive or sealant, wherein one or more of the reactive constituents is encapsulated so as to isolate it from the other constituents. These adhesive and sealant compositions may be of a number of different types: some of the more typical being those based on epoxies, urethanes, unsaturated polyesters, alkyds, and (meth)acrylates, as monomers, pre-polymers, and low molecular weight polymers or combinations thereof. Such adhesive and sealant compositions are more typically found in the form of and characterized as two or more part systems wherein the parts are combined by the applicator at the time of use. However, with the advent of encapsulation technology, it is possible to encapsulate one or more of the reactive constituents, isolating it from the other constituents, so as to produce storage stable, one-part adhesive and sealant compositions. For example, the curatives or curing and/or cross-linking agents and/or other reactants may be encapsulated and said capsules dispersed in the liquid polymerizable monomer which forms the matrix of the adhesive or sealant. Alternatively, the liquid polymerizable component may be encapsulated and the curative or curing agent adhered to the outer wall of the capsule or the liquid polymerizable component may itself be encapsulated as well.\nOftentimes these one-part adhesive and sealant compositions encompass several different microcapsules, each containing one of the curatives or curing and/or cross-linking agents and/or other co-reactive constituents alone or together with other constituents of the adhesive or sealant composition. For example, with a free-radical polymerizable adhesive or sealant composition, so long as the oxidizing agent (typically the peroxide) and the reducing agent (typically an amine and/or metallocene) are in separate capsules, the system is stable. Here two different microcapsules may be employed wherein each contains a portion of the polymerizable component and one or more of the aforementioned curing agents and/or co-reactants. Similarly, for co-reactive polymerizable systems, each of the co-reactive species is encapsulated in separate microcapsules.\nIn the case of pre-applied adhesive and sealant compositions, typically a liquid adhesive or sealant composition containing the encapsulated component(s) dispersed therein is applied to the substrate and a polymer film formed over and encasing the liquid adhesive or sealant on the substrate to which it is applied. The polymer film holds the liquid adhesive or sealant composition in place and forms a protective barrier as well as a dry to the touch surface for the adhesive or sealant composition. Here, however, the curable or polymerizable component of the adhesive or sealant is not present in an encapsulated form, i.e., is not in the form of discrete microcapsules.\nA second embodiment for the pre-applied adhesive and sealants is that where one or more of the components of the adhesive or sealant, including, in particular, the liquid polymerizable components, is encapsulated in a plurality of microcapsules which are dispersed in a liquid curable, polymerizable, or hardenable binder system. In this embodiment, one or more solid curative or curing agents or one or more solid activators, catalysts, initiators, accelerators, and the like for effecting curing or polymerization of the encapsulated liquid polymerizable component may be dispersed in the binder, without encapsulation. Alternatively, all of the components of the adhesive or sealant composition may be encapsulated in a plurality of different microcapsules which are dispersed in the liquid binder. In use, each of these modified binder systems is then applied to the intended substrate and allowed to cure, polymerize or harden; thereby binding the microcapsules to the substrate surface. Suitable binder systems may or may not co-react with the adhesive or sealant composition during cure or polymerization of the same. Most often the binders do not co-react, but instead act as a filler of the polymerizable or curable compositions.\nActivation of these encapsulated adhesive and sealant compositions is effectuated by crushing the microcapsules so that the liquid polymerizable component comes into contact with the curative, curing agent, activator, catalyst, initiator, accelerator or the like. Mixing of the ingredients is reliant upon flow caused by the application of pressure and/or relative movement of the substrates to be bonded. Thus, because of the limited mixing, it is important that such systems be as fluid as possible to maximize the opportunity for mixing. Higher viscosity materials will result in less mixing and, thus, only partial curing of the curable materials. More importantly, higher viscosity materials are more difficult, if not impossible, to encapsulate. Where encapsulation is possible, particle size becomes an issue. Typically, in encapsulating high viscosity liquids, one tends to get a large number of large and small particles or microcapsules rather than a parabolic distribution where large and small particles are few.\nWhile the advent of such encapsulated adhesive and sealant compositions has greatly broadened and/or simplified the end-use applications to which such adhesive and sealant systems are employed, they are not without remaining shortfalls. Premature fracturing of the capsules is perhaps the most prevalent of issues, especially in situations where the adhesive or sealant is used in rapid industrial manufacturing processes requiring quick delivery and application of the adhesive or sealant and, further, especially in the case of pre-applied adhesives or sealants, where the substrates to which the adhesives or sealants are pre-applied are likewise subject to such manufacturing processes, repeated handling or movement, etc. Concern for premature fracturing is not, however, limited to the application and use of the encapsulated adhesive. Such is also a great concern in the storage and handling of the microcapsules as well as the manufacture and processing of the adhesive or sealant itself, especially during incorporation of the encapsulated component into the matrix of the adhesive or sealant, or, in the case of pre-applied adhesive, incorporation of the encapsulated adhesive systems into the binder materials. Concerns with premature fracturing are heightened with systems that are thixotropic, have a high viscosity and/or incorporate fillers, especially granular fillers and fillers having sharp edges, or are subject to high shear mixing and dispensing processes.\nOftentimes to overcome concerns with premature fracturing, the thickness of the microcapsule walls is increased. This is especially found with pre-applied adhesives and sealants, particularly thread-locking or thread-sealing pre-applied adhesives and sealants, respectively. Consequently, cell walls comprising as much as 30 weight percent or more, more typically from 10% to 20% by weight, of the encapsulated component are not unusual in those one-part liquid systems to be applied at the time of use or from 15% to 25%, by weight in the case of pre-applied adhesives. However, as thicker and thicker walls are employed, less and less curable material is available with the same weight or volume of adhesive or sealant. Consequently, there is less curable material at the bond-line for forming the adhesive bond or seal. Because of the limited amount of curable material, there is poor flow and wetting of the substrate surfaces. Additionally, the large amount of shell wall fragments creates a gap between two substrates to be mated, much like grains of sand between plates of glass, which gap may be difficult to fill with the liquid curable material: again due to its limited quantity and low viscosity.\nFurthermore, as the wall thickness increases, it becomes more and more difficult to break the cell walls when fracturing is desired. Thus, greater forces are needed to ensure the release of a comparable amount of the liquid curable or polymerizable component. Because rupture of the capsules in these instances is typically as a result of finger or hand pressure, pinch rollers and/or threading of threaded elements, the degree of rupture of the microcapsules lessens as the thickness of the capsule walls increases, especially if there is not a concurrent increase in the amount of pressure applied. This is particularly so for applications other than threading applications where multiple rotations are employed. Consequently, depending upon the specific end-use application, there may be poor or low sealant or bond strength due to the release of insufficient liquid curable or polymerizable components; thus, limiting performance or even the suitability of such materials for a given application. Similarly, as oftentimes found with the bonding and/or sealing of threaded elements, a larger volume of the pre-applied adhesive or sealant is necessary in order to attain the same bond or seal volume of cured material arising from the use of liquid, as opposed to pre-applied, adhesive and sealant compositions.\nIn order to moderate the need for thicker and thicker capsule walls, it is also known to incorporate hollow microspheres as “spacer” particles in the adhesive or sealant composition. For example, Hinterwaldner (U.S. Pat. No. 4,362,566) employs hollow microspheres to enhance storage stability of the adhesive formula as well as enhance fracturing of the microcapsules during activation and application of the adhesive material. However, the addition of such microspheres adds yet another additive to the system and, because the volume of the microspheres must be accommodated, requires the use of larger volumes of the adhesive and sealant to achieve the same bond or seal volume.\nConsequently, it would be desirable to achieve encapsulated adhesive and sealant compositions in which the cell wall of the capsules is minimized and the quantity of liquid curable component is increased.\nIt would also be desirable to achieve encapsulated adhesive and sealant compositions having controlled flow characteristics, especially high or higher viscosities than are traditionally found with encapsulated adhesive and sealant compositions, and, in particular, having a fairly constant rheology over a larger temperature range, so as to enable their use in gap situations and other situations where squeeze-out or flow-out of the liquid curable material is undesirable.\nIt would also be desirable to achieve encapsulated adhesive and sealant compositions wherein, even if premature rupturing or fracturing of the capsules occurred, the composition was such that minimal, if any, reaction or polymerization of the curable or polymerizable components was able to occur.\nIt would also be desirable, especially for certain applications, to achieve the foregoing without the use of constituents or materials that are not relevant to the bonding and/or sealing of the substrates and surfaces to which they are to be applied.\nIt would also be desirable to achieve one-part storage stable adhesive and sealant compositions which are especially suited for use in high speed, industrial manufacturing processes.\nIt would also be desirable to achieve pre-applied adhesives which are able to withstand large pressures and forces without activation, especially in a pre-applied state on a given substrate."}
{"text": "Prostheses of substantially the above-mentioned construction for artrodhesis already exist in many embodiments.\nExamples of such embodiments are found in e.g. US 2010/0130978 A1. One of the drawbacks of this prior art construction is that it is adapted solely for artrodhesis (primary artrodhesis). Should there be a reason for replacing the prosthesis in order to again make the joint flexible, the entire prosthesis must be replaced. Furthermore, the prior art prosthesis does not include any member of the previously used prosthesis for making the joint flexible (artroplasty). Another drawback is the limited possibility to adjust the setting of the prosthesis members relative to each other and that adjustment must be performed by means of special tools.\nOther similar constructions for artrodhesis of e.g. knee joints are found in DE 3347055 A1 and US 2009/0171463 A1.\nIn SE 528545 C2, a prosthesis for replacement of a joint (artroplasty) is described. This prior art prosthesis is configured substantially as defined above. This prior art construction however, is not adapted for artrodhesis and members thereof can not be used therefor."}
{"text": "1. Field of the Invention\nThe present invention relates to electronic semiconductor devices, and, more particularly, to erasable programmable memory devices and methods of fabrication.\n2. Description of the Related Art\nNonvolatile semiconductor memory devices based on metal-oxide-semiconductor field effect transistors (MOSFETs) were first proposed in 1967 (see Sze, Physics of Semiconductor Devices, pp. 496-506 (Wiley-Interscience 2d Ed 1981)). These devices store a bit of information as the presence or absence of a quantity of electrical charge on a floating gate that is located so that the charge affects the threshold voltage of a MOSFET. Currently, MOSFET nonvolatile memory devices include EPROMs, EEPROMs, and Flash EEPROMs. EPROM (erasable programmable read only memories) cells store a bit of information as a quantity of electrons on a floating gate; the electrons are avalanche injected into the floating gate from the drain-end of the cell transistor channel and are erased for all cells simultaneously by photoemission under ultraviolet irradiation. The Floating gate Avalanche injection MOS transistor is typically called a \"FAMOS\" transistor. The density of EPROMs had reached 1 megabit by 1984, but the inconvenience of the ultraviolet erase has led to the development of EEPROMs and Flash EEPROMs.\nEEPROMs (electrically erasable programmable read only memories) both program and erase information for a single memory cell by tunneling charges through insulators: the FLOTOX version uses a floating gate as with an EPROM but with a thin tunneling oxide between the floating gate and the transistor drain; the textured poly version uses tunneling between three polysilicon gates with one of the gates a floating gate; and the MNOS version uses a stacked oxide and nitride gate insulators and stores charges in traps in the nitride and programs and erases by tunneling through the oxide to the channel region. See generally S.Lai et al, \"Comparison and Trends in Today's Dominant E.sup.2 Technologies,\" Digest of Technical Papers, pp. 580-583 1986 IEDM.\nFlash EEPROMs are hybrids that program in the manner of either EPROMs (avalanche injection) or EEPROMs (tunneling) and erase in the manner of EEPROMs (tunneling) but with the erasure generally limited to bulk electrical erasure of the entire memory analogous to the ultraviolet light erasure of an EPROM.\nThe trend to larger scale integration has demanded small, highly packed memory cells with low power dissipation, and the desire for flexibility of multiple reprogramming has demanded durable tunneling oxides. Consequently, many variations of the basic EPROM, EEPROM, and Flash EEPROM cells have been introduced. For example, McElroy, U.S. Pat. No. 4,373,248 shows EPROM cells in an array with a set of parallel continuous buried n.sup.+ diffusion lines in a silicon substrate acting as sources and drains (bitlines) for the floating gate cells with a second set of parallel continuous polysilicon lines over the floating gates acting as control gates (wordlines); the polysilicon lines are perpendicular to the buried diffusion lines. J. Esquivel et al, High Density Contactless, Self Aligned EPROM Cell Array Technology, 1986 IEDM Digest of Technical Papers, pp. 592-595, Mitchell, U.S. Pat. No. 4,597,060, and Sugiura et al, U.S. Pat. No. 4,451,904 also have buried bitlines and perpendicular polysilicon wordlines in an EPROM array.\nMiyamoto, U.S. Pat. No. 4,642,673 has a floating gate EEPROM cell array with parallel continuous buried diffusion lines acting as control gates and each cell has a floating gate transistor plus a select transistor in series with the floating gate extending over a diffusion line; the source/drain contacts are to metal lines (bitlines and source lines) perpendicular to the diffusion lines, and gates for the select transistors are polysilicon lines (wordlines) parallel to the diffusion lines. Note that the series select transistor is added to avoid the case of an over-erased floating gate transistor going into depletion mode and conducting even with no voltage applied to the control gate.\nD. Guterman, U.S. Pat. No. 4,590,504 has an EEPROM array with each cell having a buried source region connected to a buried ground line and a buried drain region including a remote drain portion for tunneling into the floating gate connected to a metal bitline plus parallel polysilicon control gate line and select transistor gate line.\nF. Masuoka et al, A 256K Flash EEPROM Using Triple Polysilicon Technology, 1985 ISSCC Digest of Technical Papers, pp. 168-169, has a Flash EEPROM cell array with first level polysilicon erase lines, second level polysilicon floating gates, and third level control gate lines (wordlines); the floating gate transistor and the select transistor are merged by having the floating gate over only a portion of the channel region. The transistor source connects to a diffusion line and the drain connects to a metal line over the polysilicon lines; the diffusion lines and the wordlines are parallel, and the erase lines and the metal lines are parallel and are perpendicular to the diffusion lines and wordlines. The floating gate is programmed by avalanche injection as in an EPROM, and erasure of the floating gate is by tunneling to the adjacent erase line.\nHowever, the known EEPROMs and Flash EEPROMs have problems of large cell size limiting packing density, isolation and programmability for dense arrays of cells and complex processing."}
{"text": "This disclosure relates to test and measurement instruments, and in particular, to an acquisition apparatus of a test and measurement instrument having an arbitrary multiband overlay mixer apparatus and associated techniques for bandwidth multiplication.\nThe performance of test and measurement instruments, including digital oscilloscopes, is limited by the input bandwidth of the acquisition system. Conventional components for acquiring high frequency signals include sophisticated track and hold sampling circuitry, analog-to-digital conversion circuitry, high speed memory, and the like. However, traditional track and hold sampling and analog-to-digital circuitry require large investments of time and capital to effectively build and deploy.\nOther approaches have involved the use of mixer technology in front of a digitizer platform to allow the system to have wider bandwidth, and a lower cost in development effort, compared to a standard digitizer platform architecture. In U.S. Patent Application Publication 2004/0128076 to Pupalaikis, et al., a real-time oscilloscope is disclosed with an increased usable bandwidth. The real-time oscilloscope splits the input signal into multiple split signals. One split signal is digitized. Simultaneously, the other split signals are frequency shifted to a baseband frequency range and digitized. The digitized frequency-shifted signals are frequency shifted to their original frequency range and then combined with the other digitized signals to create a representation of the input signal. By frequency shifting sub-bands of the input signal to be within the bandwidth of their respective digitizers, an input signal having a frequency range larger than the input bandwidth of a digitizer may be acquired using the lower bandwidth digitizers.\nHowever, a need remains for a multi-band overlay apparatus and method that allows for arbitrary mixer local oscillator coefficients having an inverse matrix relationship to local oscillator coefficients used in the reconstruction of an input signal. In addition, it would be desirable to simultaneously digitize two wide band input signals using four digitizers of narrower bandwidth in a particular arrangement that provides for the use of arbitrary coefficients and an efficient reconstruction algorithm."}
{"text": "The present invention relates to long term or semipermanent storage of spent nuclear fuel assemblies, and more particularly to a storage rack for spent fuel assemblies and a method of storage of such fuel assemblies to maximize the capacity of the rack while preventing criticality.\nNuclear reactors consist of an array of fuel rods containing the nuclear fuel. The fuel rods are metal tubes, typically from 8 to 15 feet in length and about 1/2 inch in diameter, and are supported in groups in fuel assemblies which may comprise a considerable number of rods. The large reactors utilized for power generation contain a large number of these fuel assemblies arranged in a suitable configuration.\nAfter an extended period of operation, the irradiated or spent fuel assemblies must be removed from the reactor and replaced. The spent fuel rods contain residual amounts of the original fuel material, and varying amounts of numerous fission products resulting from fission of the nuclei of the original fuel and other nuclear reactions, as well as from radioactive decay of initially-formed fission products. Certain of these materials are themselves fissionable. Many of the fission products are highly radioactive, at least initially, and thus produce considerable heat, while the entire fuel assembly is dangerously radioactive. The fuel rods can be reprocessed by chemically separating the fissionable material for reuse as fuel and recovering various other fission products, such as certain rare earth elements, for example, which have substantial commercial value.\nSuitable facilities must be provided for the storage of these highly radioactive fuel assemblies after removal from the reactor until they can be reprocessed or otherwise disposed of. Such storage presents serious problems since the fuel assemblies are initially highly radioactive and generate a great deal of heat. They must, therefore, be kept submerged in water which serves as a coolant to prevent overheating as well as a radiation shield and moderator for the fast neutrons which are still being emitted. It is also necessary to be sure that the assemblies are stored in a manner that will prevent criticality of the collection of fuel assemblies while keeping the space required to a minimum. After some period of time, the heat generated and the radioactivity of the fuel assemblies decline, since many of the fission products have relatively short halflives, and the nature of the storage problem changes as both the heat to be dissipated and the radiation hazard decrease.\nIn our prior U.S. Pat. No. 4,010,375, there is disclosed a storage rack for spent nuclear fuel assemblies which is primarily intended for temporary storage of fuel assemblies in a water-filled pit. The rack consists of a checkerboard array of storage cells with the spent fuel assemblies placed in alternate cells. The intervening cells are filled with water which functions as a moderator and as a coolant, and also include a poison or neutron-absorbing material. This arrangement prevents criticality of the collection of fuel assemblies while maximizing the capacity of the rack for relatively short term storage. We have also proposed, in our copending application Ser. No. 851,038, filed Nov. 14, 1977, to use a similar rack structure for very long term or permanent storage, and to maximize the capacity by completely filling all the cells in the array after sufficient irradiation of the fuel and subsequent radioactive decay has occurred to make this safely possible. Concrete shielding may be provided for permanent storage. There is also a need, however, for relatively long term or semipermanent storage of spent nuclear fuel assemblies in a manner which will permit the storage of a maximum number of such assemblies in a given space, with complete safety, for a relatively long period until they can be disposed of by reprocessing, placing them in permanent storage, or otherwise.\nIn the conventional design of spent fuel storage facilities, criticality is prevented by means of the spacing, or pitch, between adjacent fuel assemblies in the storage rack. There is a possibility that it may at some time be necessary to completely unload the reactor to make repairs or inspections inside the reactor pressure vessel, and it is, therefore, assumed that it may be necessary to place nearly fresh or unirradiated fuel assemblies in the storage rack during such repairs or inspection. In conventional designs, therefore, the pitch between adjacent fuel storage locations has been determined on the basis of the reactivity of fresh or unirradiated fuel assemblies. This, of course, requires a much larger pitch than is needed to prevent criticality with spent fuel assemblies after discharge from the reactor, and results in reduced storage capacity of a given space. Furthermore, once the pitch has been determined in a conventional design, it cannot be changed without rebuilding the storage rack. The arrangement and dimensions of such a rack are such that additional spent fuel assemblies cannot be placed between the initial storage positions both because the usual designs do not provide for accommodating such additional fuel assemblies and because the dimensions are generally too small."}
{"text": "A wide variety of techniques are used to manufacture prostheses and models that have the feel and appearance of living tissue. Methods for manufacturing simulated skin which have the texture and appearance of animal or human skin are known in the art. However, models of animals or prostheses prepared from artificial skin often lack the feel of living tissue because the artificial skin is not connected to underlying structures in a manner that simulates the characteristics of connective tissue."}
{"text": "The invention relates to a surgical punch with a tubular shaft, the distal end of which forms an annular cutting edge, a cutting member which is mounted for longitudinal displacement in relation to the shaft, interacts with the cutting edge and is held on an advancing element mounted for displacement in the tube, and an actuating member for displacing the advancing element in the shaft, this member being arranged at the proximal end of the shaft.\nSurgical punches of this type are used to introduce openings into vascular walls, for example for attaching a bypass vessel. In this respect, a cutting member which can, for example, be of a conical design is brought close to a cutting edge with its circumferential edge, wherein the vascular wall is entrained between cutting member and cutting edge. As a result, a cut is made in the vascular wall along the cutting edge, i.e. circular openings can, for example, be introduced into vascular walls.\nIn the case of conventional instruments, a different punch has to be used for each opening size; moreover, these instruments can be used only for such a time as the cutting edge of the cutting member generates a clean cut, i.e. the cutting edge has the desired sharpness and is not damaged in any way. If the cutting edge becomes blunt or damaged, the instrument has to be replaced."}
{"text": "With the development of touch technology, touch input devices have been widely used in many end electronic products. Because being capable of providing a better controlling and various interface operation ways, touch input devices have replaced the traditional input devices, such as mouse and keyboard, and accordingly the end electronic products need less space and consequentially solve the issue of being inconvenient to carry the peripheral input device.\nBased on the touch technology, touch panels can be categorized to resistive, capacitive, optical and sound wave types; wherein the capacitive touch panel currently is relatively more popular in the market and has a higher market share. However, there still has some problems needed to be solved in the conventional capacitive touch sensing system. For example, a capacitive touch sensing system may result in an error operation while a touch panel thereof is being touched by a palm or a wrist of a user; or, for eliminating the effect of wrong touching, the sensing system may be designed to have a complex operation, which may consequentially lead to a long operation time."}
{"text": "1. Field of the Invention\nThe invention relates generally to presses and, more particularly, to shell presses and associated methods for forming container closures or ends, commonly referred to as shells. The invention also relates to die assemblies for shell presses.\n2. Background Information\nThe forming of can ends or shells for can bodies, namely aluminum or steel cans, is generally well-known in the art.\nThere is an ongoing desire in the can-making industry to manufacture shells as rapidly and efficiently as possible. Among the ways companies have attempted to achieve these objectives are: (1) to increase the number of pockets in the die set, within which shells can be formed; and (2) to increase the speed (e.g., strokes per minute (spm)) at which the shell press operates. In general, with each stroke of the shell press ram, one shell is formed in each tooling pocket of the die assembly. Thus, a 24-out die assembly, for example, which has 24 tooling pockets, is capable of forming 24 shells, per stroke. U.S. Pat. No. 5,491,995, which is hereby incorporated herein by reference, discloses an example of a relatively high capacity (e.g., without limitation, operating speed of up to 400 spm, or more) end shell manufacturing system having a 24-out die assembly.\nHowever, forming shells at relative high speeds generates heat. The heat, which is caused by the friction associated with drawing the metal over forming surfaces of the die assembly and/or clamping the metal between various pressure pads and drawing it through reduced tooling clearances to provide a desired shape, can be excessive, resulting in thermal expansion of the die shoes. Among other disadvantages, such thermal expansion undesirably shifts tooling and/or reduces critical clearances between cutting and/or forming tools. Consequently, tooling wear or damage can result and/or certain features of the end shells are manufactured out-of-specification. For example, thinned spots can be created in the material from which the end shell is manufactured, leading to a loss in buckle pressure performance in the final product.\nThe foregoing difficulties have been exacerbated by the development of new shell designs having aggressive material thicknesses and shapes. For example, some shells require reduced material thickness and/or have a relatively complex geometry. Such shapes often necessitate additional pressure pads and increased forming pressures in order to properly manufacture the end shells.\nPrior proposals that attempted to address thermal expansion of the die assembly tooling (e.g., without limitation, upper and lower die shoes) involved aligning the upper tooling with respect to the lower tooling in the die assembly in a manner intended to compensate for the thermal expansion. Other proposals require coolant (e.g., chilled water) to be pumped throughout the die assembly, for example, to reduce the rate and amount of thermal expansion of the die shoes. However, estimating and establishing the proper aligning of the upper tooling with respect to the lower tooling is a time-consuming process, and it can be difficult to maintain the desired alignment. Similarly, systems that add coolant or other suitable additional cooling or heating mechanisms to the die assembly to compensate for thermal expansion, are costly to install and maintain.\nThere is, therefore, room for improvement in shell presses, and in die assemblies and associated methods therefor."}
{"text": "When applying water based adhesives by hand spray techniques or automated/machine controlled spray techniques for assembly of cushioning materials, such as for the furniture and bedding industries, there is a problem with adhesive overspray. The overspray presents itself as a fog in the factory that can carry long distances from the actual application area of the factory. This fog also creates a nuisance dust health hazard for the employees. Lastly, the fog or overspray is wasteful of resources as the adhesive is lost and not used for its intended purpose. This overspray not only gets onto the employees that apply the adhesive, but also contaminates nearby equipment, finished product, or raw materials in inventory, and contaminates air conditioners, heaters, and lighting.\nOne solution has been to set up air extraction hoods in the spray area. This works relatively well when the filters are maintained and the types of parts that are being assembled are small. However when making larger items such as mattresses, large sofa cushions, and the like, the usefulness of an air extraction hoods is negated because of the impracticality of extraction hoods that are sized for large items.\nAlso there have been attempts to control the overspray fog by using low fogging air atomized guns such as the DUX or EasyFlow Laminair spray gun. Although these spray devices minimize the over spray when adjusted properly, they are dependent on the spray operators not adjusting the settings as they can easily be misadjusted and create fog.\nAnother solution has been to use different types of adhesive bases other than water base. Solvent based adhesives and hot melt adhesives when sprayed do not create a fog. These types of adhesives work well to eliminate the overspray but have other problems.\nSolvent based adhesives contain hazardous materials and often are flammable. They require air extraction equipment to reduce the flammability hazard as well as the health hazards to employees. Also solvent adhesives do not adhere some types of visco-elastic foams.\nHot melt adhesives typically do not bond foam cushion substrates as well as water based or solvent based products. Hot melts also require melt tanks and heated hose and this equipment is more expensive on a per gun basis than water based or solvent adhesives.\nAnother solution is the roll coating of water based adhesive rather than spray application. Roll coating eliminates the over spray, but suffers additional problems because the rollers are exposed to the atmosphere. As such, during any down time at all, the adhesive on the rollers can coagulate, causing inconsistent application of the adhesive. In addition, at the end of a shift, the workers must clean the rollers, adding to the system downtime and taking away working time from the workers. Further still, rollers do not allow a control of the application rate over a surface. Although roll coating provides a consistent application of adhesive across an entire surface, sometimes it is advantageous to vary the application rate of the adhesive. For example, it may be advantageous to use more adhesive in one area, and less in another, thereby using less adhesive overall."}
{"text": "In a general digital broadcasting system, because data are transmitted through a limited frequency bandwidth, a time dispersion effect, which means that pulse energy of a predetermined symbol is dispersed to neighboring symbols, generates interference to the neighboring symbols. Furthermore, the transmitted data is affected from various channel distortions.\nAs channel distortion, there are such as a multipath, a frequency offset and a phase jitter. The channel distortion occurs an Inter-symbol Interference (ISI), which means interference or a distortion between a predetermined symbol and its neighboring symbols and generates a problem to receive desired data in the digital transmission system.\nTo prevent the channel distortion and decrease a symbol error due to an ISI, the general receiver, e.g., a digital broadcasting receiver, uses a channel equalizer.\nMost of the communication channels need an adaptive equalizer for updating a tap coefficient according to time because the distortions in the general digital broadcasting system are time variant.\nConfiguration of a conventional channel equalizer is described in FIG. 1.\nAs shown in FIG. 1, the conventional channel equalizer includes a digital filter 11, a symbol detector 12, a tap coefficient updating unit 13, a training sequence storage 14, a statistical data calculator 15, a switch 16 and an equalizer input signal storage 17.\nThe digital filter 11 eliminates an inter symbol-interference (ISI) causing a distortion of a baseband signal received to a digital broadcasting receiver. The symbol detector 12 receives the output signal of the digital filter 11 and calculates a decided data by comparing the output signal to a predetermined threshold to thereby output the decided data.\nThe tap coefficient updating unit 13 receives output signals of the equalizer input signal storage 17 and the digital filter 11 and an error data selected by the switch 16 to thereby update a tap coefficient of the digital filter 11.\nAlso, the training sequence storage 14 stores the training data sequence already known in a transmitter of a digital broadcasting system. The training data sequence is read in a training mode and is outputted to the tap coefficient update unit 13.\nThe statistical data calculator 15 calculates a statistical error in order to output the statistical error to the tap coefficient update unit 13 in a blind mode.\nThe switch 16 selects one of the training storage 14, the statistical data calculator 15 and the symbol detector 12 according to an operating mode and outputs error data outputted from the selected one to the tap coefficient update unit 13. Then, the tap coefficient update unit 13 calculates an error signal corresponding to the error data and updates the tap coefficient by using data corresponding to the tap coefficient of the digital filter 11 in the equalizer input signal storage 17 to output the tap coefficient to the digital filter 11.\nAs the channel equalizer, the decision feedback equalizer (DFE) is used broadly in the digital broadcasting receiver. In general, the an eye diagram of an output signal in the decision feedback equalizer is opened, wherein the eye diagram is a kind of factor for determining performance of the equalizer, i.e., a function of making an output signal decision rightly and easily. Also, if an output signal of the symbol detector is determined correctly, the feedback filtering block deletes an Inter Symbol Interference (ISI) according to the predetermined symbol and there is not a problem such as a noise amplifying phenomenon in channel equalization of the a linear equalizer. Therefore, the decision feedback equalizer is widely used in the digital broadcasting receiver.\nTherefore, for an appropriate use, it is important that the output signal of a symbol detector should not have a decision error and, for the most of all, the eye diagram is opened.\nFor the eye diagram is opened in Advanced Television System Committee (ATSC) digital broadcasting system which is American standard for a terrestrial digital television, it is used a method for opening the eye diagram by inserting one training sequence segment, wherein one segment is 208 bytes, in every 312 data segment although a data efficiency is decreased.\nHowever, since the inserted training sequence segment is short under a multipath environment having a long ghost, it is often failed to open the eye pattern. In particular, though there is the training sequence, it is often failed if a tap coefficient is not enough long to converge or if the environment is under an inferior multipath environment having a time-varied channel, a long ghost or a large level signal ghost. If the eye diagram is not opened, possibility of the decision error in the symbol detector is very high. As a result, it can be generated an error propagation problem that the decision error is accumulated through a feedback loop of the decision feedback equalizer.\nTherefore, it is required a method for decreasing the decision error during data period having no training sequence segment; and, in particular, total tap energy should be decreased in order to decrease the decision error by opening the eye diagram of the decision feedback equalizer.\nAt first, most of conventional methods for decreasing the decision error use a viterbi decoder having a decoding delay. There is a method for making a predetermined delay in an equalizer tap coefficient adjustor be identical to the decoding delay of the viterbi decoder. Herein, this method is proposed by G. Long, entitled “The LMS Algorithm with Delayed Coefficient Adaptation”, IEEE Trans. Acoust., Speeach, Signal Processing, vol. ASSP-37, October 1989.\nIn addition, there is a method for solving a decoding delay of the viterbi decoder by periodically adding an interleaver and deinterleaver, proposed by M. V. Eyuboglu, entitled “Detection of Coded Modulation Signals on Linear, Severely Distorted Channels Using Decision-Feedback Noise Prediction with Interleaving”, IEEE Trans. Commun., vol. COM-36, pp. 401-409, April 1988 and in U.S. Pat. No. 4,833,693 issued to Eyuboglu, May 23, 1989.\nBecause above described methods for decreasing the decision error are used in the viterbi decoder having the decoding delay as TDB-1 behind the equalizer served as the symbol detector of the decision feedback equalizer, it is required that additional devices for deleting the decoding delay. Also, for the viterbi decoder behind the equalizer has enough performance, the TBD should be 5 times more than the number of memories in the trellis encoder used for encoding process. However, the decoding delay should be as small as possible in order to use the output of the viterbi decoder as a feedback input of the decision feedback equalizer.\nParticularly, the decoding delay is not TBD-1 but 12×(TBD-1) in the digital broadcasting system which uses 12 TCM encoders by the trellis code interleaver as shown in FIG. 5; and, approximately, the decoding delay becomes 168 because there are 2 memories in the TCM encoder. It is very inefficient that the viterbi decoder having 168 decoding delay is actually used in the digital broadcasting system.\nTherefore, for using the viterbi decoder as the symbol detector of the decision feedback equalizer in the digital broadcasting receiver, the decoding delay is as small as possible and, further, it is the best that there is no decoding delay. And a complexity for decoder implementation should be as low as possible.\nMeanwhile, in order to decrease the tap energy of the feedback filter, there are methods for increasing the tap number which removes post ghost and for changing channel property of a receiving signal by using a beam-forming or a channel-matched filter.\nThe method for increasing the tap number of the feed forward filter is inefficient and the improvement of the performance is small in comparison to the increased tap number. The method for changing the channel property by using the channel-matched filter, which is proposed in a document by Richard Citta, entitled “A VSB Receiver Designed for Indoor and Distributed Transmission Environments”, IEEE 52nd Annual Broadcast Symposium, Oct. 9-11, 2002, is more effective than the method for increasing the tap number.\nBecause the channel equalization method proposed by Richard Citta generates the channel-matched filter based on over-sampling data and uses a fractionally-spaced equalizer, the complexity is very high. Also, since a simple slicer is used as the symbol detector, the error propagation problem due to the decision error can occur.\nTherefore, development for a channel-matched filter having low complexity on a symbol basis and a symbol detector having rare decision error are highly required.\nDisclosure\nTechnical Problem\nIt is, therefore, an object of the present invention to provide an apparatus and a method of decision feedback equalization that make a channel property of an inferior receiving signal to mild by using a channel-matched filter and decreases decision errors of symbol detector output signals by using a trellis decoder with decreased complexity, whose trace back depth is 1 (TBD=1).\nThe other object of the present invention provide a symbol detection method of the trellis decoder with decreased complexity, of which TDB is 1, which can use an output signal of trellis decoder as an input signal of the feedback filter of the decision feedback equalizer by decreasing a decoding delay.\nTechnical Solution\nIn accordance with one aspect of the present invention, there is provided a decision feedback equalizer in a terrestrial digital broadcasting receiver, including: a channel estimating unit for estimating a channel of a symbol-based receiving signal based on the receiving signal and a training sequence; a channel-matched filtering unit for changing a channel property of the receiving signal by maximizing a signal-to-noise ratio (SNR) of the estimated channel; an input signal storing unit for storing a receiving symbol of which channel property is changed by the channel-matched filtering unit; a channel equalizing unit for performing a decision feedback equalization by repeatedly filtering the receiving signal which passed through the channel-matched filtering unit; a trellis decoding unit for detecting a symbol, which is decision data, based on trellis decoding algorithm with decreased complexity, whose trace back depth is 1, from channel equalized receiving symbols and outputting the symbol in a decision directed mode; a statistical data calculating unit for calculating statistical error data used in a blind mode and outputting the statistical error data; a training sequence storing unit for storing the training sequence; a switching unit for selecting a mode among the training mode, the decision directed mode and the blind mode; an error signal calculating unit for calculating an error signal by comparing an output signal in the mode selected by the switching unit to an output signal of the channel equalizing unit; and a tap coefficient updating unit for updating a tap coefficient to be provided to the channel-matched equalizing unit based on the error signal, the output signal of the trellis decoding unit and the output signal of the input signal storing unit.\nIn accordance with one aspect of the present invention, there is provided a decision feedback equalizing method in a terrestrial digital broadcasting receiver, including the steps of: a) estimating a channel of a symbol-based receiving signal based on a receiving signal and a training sequence; b) changing a channel property of the receiving signal in order to maximize a signal-to-noise ratio (SNR) of the estimated channel by passing the receiving signal through a channel-matched filter; c) determining a parameter used for a decision feedback of the receiving symbol whose channel property is changed and initializing a channel equalization parameter; d) detecting a symbol from an output signal of an equalizer in a specific time index signal according to the determined parameter based on a trellis decoder with decreased complexity, whose trace back depth is 1; e) calculating statistical error data used in a blind mode; f) selecting one mode among a training mode, a decision mode and the blind mode; g) calculating an error signal by comparing an output signal of the mode selected in the step f) to output signal of equalizer in a decision directed mode, and updating a tap coefficient based on the error signal; and h) performing a decision feedback equalization based on the updated tap coefficient.\nIn accordance with one aspect of the present invention, there is provided a symbol detecting method for channel equalization in a terrestrial digital broadcasting receiver, the method including the steps of: a) calculating an absolute distance pairs between symbol pairs in an input signal of a symbol detector and a trellis diagram; b) selecting an absolute distance having a small value for each absolute distance pair among absolute distance pairs; c) calculating an accumulated absolute distance by adding a previous absolute distance to a current calculated absolute distance for each state in the trellis diagram in a time index; d) deleting the accumulated absolute distances except the smallest accumulated distance for each state in the trellis diagram in the time index; e) selecting a state in which the accumulated absolute distance is smallest among all states shown in the trellis diagram in the time index and obtaining an output signal of the symbol detector, which is the trellis decoding unit, from a branch shown in the trellis diagram transited to the selected state; and f) repeatedly performing the steps a) to e) for each symbol time index.\nThe present invention relates to an apparatus and a method for decision feedback equalization having a channel-matched filter and a trellis decoder in a terrestrial digital broadcasting receiver, wherein the channel estimator estimates the channel of the receiving signal based on the receiving signal and the training sequence, generates the channel-matched filter based on information of the estimated channel, equalizes the receiving signal passed though the channel-matched filter by using the decision feedback equalizer having the viterbi decoder with decreased complexity, whose TBD is 1, and thereby the channel equalization is performed effectively under an inferior environment such as in a room or in mobile\nThe present invention can change the channel property of the receiving signal under an inferior environment such as in a room or in mobile to be mild by placing the channel-matched filter generated from the symbol-based channel estimator in front of the conventional decision feedback equalizer, and decreases the decision error of the symbol detector by using the viterbi decoder with decreased complexity, whose TBD is 1, in stead of the simple slicer served as the symbol detector in the conventional decision feedback equalizer and thereby a convergence speed and stability of the digital filtering block are increased, and the residual MSE is decreased in a normal state after the convergence.\nAdvantageous Effect\nAs mentioned above, the present invention can change the channel property of the receiving signal under an inferior environment such as in a room or in mobile to be mild by placing the channel-matched filter generated from the symbol-based channel estimator in front of the conventional decision feedback equalizer, and decreases the decision error of the symbol detector the viterbi decoder with decreased complexity, whose TBD is 1, in stead of the simple slicer served as the symbol detector in the conventional decision feedback equalizer and thereby a convergence speed and stability of the digital filtering block are increased, and the residual MSE is decreased in a normal state after the convergence."}
{"text": "1. Field of the Invention\nThe invention relates to a device and methods for profiling railway rails and, more particularly, to a device and methods for vehicle wheel profile-matching railway turnouts and crossings.\n2. Description of Related Art\nThe contact rolling surface found along the length of a standard railway crossing, such as a steel frog or a diamond crossing, has been pre-defined for many years. When a vehicle passes through a typical rail crossover, the wheels pass over a raised wing surface, which is designed to support and maintain the wheel profile for a relatively level passage. Rail crossover, as used herein, means a rail turnout or crossing allowing a rail vehicle to be guided from one set of rails to another. Insert and solid steel frog designs have been supplied for many years and are used by every major railway. Very little advancement has taken place with regard to technology designed to precisely match an existing rail crossover to a specific existing vehicle wheel profile or to alternate wheel profiles. Vehicle wheel widths and tapered tread profiles are often custom designed to suite a particular transit operator's requirements. As a result, when a standard crossover design is installed, oftentimes the transfer surfaces are not continuously levelly supported through the entire length of the crossover casting. Without matching the crossover profile to the wheel requirements, level rolling discontinuities often result in dropping of the wheel tread surface, thereby, imparting high impacted forces and accelerated wear/damage to the track and vehicle bogie system. Bogie impact forces from as little as 5 mm of height have resulted in up to 30-40 G of impact force during higher vehicle operating speeds and significantly have reduced operating life expectancy for track and vehicle components.\nWhen field repair is required on a damaged crossover, rolling surface areas are built-up through a special certified welding process and the profile manually shaped back to the original (usually flat) profile with a hand grinder. The finished surface contour is usually visually inspected with a straight edge to verify that the finished product is returned as close as possible to the original rail surface contour. Without a technique to accurately apply and precisely shape the length of the crossover to a fully level supported “wheel tread matching” profile, and without an accurate means to quickly and easily verify an alternate properly supporting profile has been provided, rolling axle instability due to wheel passage and track/equipment wear will continue to occur.\nTypical technology used to shape the crossover surface after welding includes grinding by use of a hand grinder or portable grinding platforms that are either secured directly over the crossover casting or are guided to provide a flat and level profile over a relatively short linear path. However, after grinding a section, the finished surface is only visually inspected and checked with a straight edge or simple hand tools to ensure reworked sections are uniform with no surface gaps present. These methods do not apply to a revised raised wing or point sections to better support the transfer of the wheel path. Therefore, wheel load impacts are not prevented, only maintained at a lower level after rework. This results in less accurate correction to the surface profiles of the crossovers with limited verification. Often, surface discontinuities are still present after modification resulting again with initially lower impact forces. The previously used methods included minor modifications, which failed to correct the source of the impact forces.\nThese high impact forces have often led to the reduction of operating speeds through crossovers and continuous repair and maintenance for both vehicles and rails. The cost and time required to routinely repair rail crossovers is an ongoing concern.\nThere exists a need for a precision and efficient method and device to greatly eliminate impact forces due to imperfectly matched rail crossover profiles and vehicle wheel profiles."}
{"text": "The performance of a gas turbine engine is sometimes limited by the maximum temperature to which certain components can be exposed. These components, such as turbine blades and vanes, various seals and spacers, combustors, and other components, are sometimes actively cooled with air from the engine's compressor, so that the component metal temperatures can be maintained at acceptable levels while being exposed to gas temperatures that are in excess of permissible metal temperatures.\nThis active cooling comes with an inherent trade off. Although higher gas temperatures tend to improve the performance of the engine, the use of air from the compressor for active cooling tends to lower the performance of the engine. It is imperative that the scheme for active cooling use as little air as possible, so that performance gains are maximized.\nIn order to cool that various components efficiently, numerous component internal cooling flowpath designs have been used to optimally transfer heat from the component into the cooling air. Examples of such cooling flowpath designs can be found in U.S. Pat. Nos. 4,930,980 issued to North et al.; 4,236,870 issued to Hucul, Jr. et al.; 4,946,346 issued to Ito; 5,062,768 issued to Marriage; and British Patent No. 2,246,174A issued to Lings et al.\nThese flowpath designs provide cooling air to various regions of the component, and often use small features to enhance heat transfer. Examples of cooling flowpath designs with small features include U.S. Pat. Nos. 4,770,608 issued to Anderson et al.; 4,962,640 issued to Tobery; 5,281,084 issued to Noe et al.; and 5,193,975 issued to Bird et al. Small sized features such as pins, pedestals, walls, and passageways permit an increased number of such features to be included in the cooling flowpath. A large number of cooling pedestals, for example, increases the surface area of the component exposed to the cooling air. A small size for cooling passageways permits not only a larger number of passageways, but may also increase the convective heat transfer coefficient of the passageway.\nHowever, use of small features such as pins, pedestals, walls, and passageways poses a problem during manufacturing. These features are often cast or etched onto one workpiece which must then be joined to another workpiece. Joining methods often include welding, brazing, or diffusion bonding. Examples of cooling flowpath designs in which portions of the workpieces are brazed include U.S. Pat. Nos. 5,263,820 issued to Tubbs; 4,786,234 issued to Readnour; 4,505,639 issued to Groess et al.; 4,482,295 issued to North et al.; 5,419,039, 5,392,515, and 5,405,242 issued to Auxier et al.; 5,383,766 issued to Przirembel et al.; 5,193,980 issued to Kaincz et al.; and British Patent Specifications No. 1,285,369 issued to Steel et al. and 1,299,904 to Curtiss Wright Corporation.\nBrazing of flowpaths with small geometric features presents the problem of excess brazing material interfering with the intended use of the feature. For example, excessive brazing material within a cooling air passageway can reduce or completely eliminate airflow through the passageway, with subsequent local or general overheating of the gas turbine component. It is important that excessive brazing material be kept from areas adjacent to various cooling flowpath features. U.S. Pat. No. 4,507,051 issued to Lesgourgues et al. discusses brazing of actively cooled gas turbine components.\nOne manner of depositing brazing material is by electrophoresis. U.S. Pat. Nos. 5,336,382 issued to Bodin, and 5,411,582 issued to Bodine, describe use of electrophoresis with a mucilage remarkable for its high viscosity, the muciliage containing brazing powder with grains as large as 53 microns. Outer areas in which brazing material is not desired are masked by a suitable adhesive masking tape, and inner cavities are masked with silicone plugs.\nWhat is needed is a method of brazing components in which brazing material can be accurately deposited near small features. The present invention provides this method in a novel and unobvious way."}
{"text": "1. Field of the Invention\nThe present invention pertains to an eyeglass end face machining method, particularly to the polishing to a mirror polishing that is performed on the end face after bevel edging, or the planing, such as machining to a mirror polishing, that is performed on the end face after edging.\n2. Description of the Related Art\nThe lens end face of rimless eyeglasses lenses usually referred to as three-piece eyeglass lenses is exposed and not covered by a rim, etc., and therefore, they must have a surface that has been polished until glossy. In response to this need, technology has been presented whereby eyeglass lenses, whose end face has thus far been smoothed manually in order to obtain a face that has been polished until glossy, are mechanically polished by placing a movement mechanism with tracing capability in the polishing wheel part (for instance, Japanese Patent Laid-Open No. Sho 64-87144). This grinds inclined faces, such as the end face of polyhedron cut lenses, etc., and although the shape around the eyeglass lens is complex because of the polyhedron cut, the end face itself, which becomes the surface to be grounded, is a flat surface and simple. Consequently, the above-mentioned technology cannot be used when the surface to be polished itself has a complex shape, such as lens end faces with a bevel. Now, because the lens end face with a bevel is usually concealed by the rim of the frame and there is no need to polish the bevel faces, a lens end face with a bevel itself is usually not polished.\nHowever, there has been a demand in recent years for thin rims in order to obtain frames that are more lightweight and fashionable, etc., and it is often the case that if the lens fitted into the rim is a strong-minus-power lens with a thick edge, the lens will protrude from the rim of the frame. It is pointed out that the bevel faces remains white when polishing of the lens end face is completed by bevel-polishing and this poses a problem aesthetically. Polishing the bevel surface that remains white until it is transparent is only accomplished by buff polishing the bevel surface by hand, etc., and this takes time and increases cost.\nThe objective of the present invention is to solve the above-mentioned problems with prior art by mechanically polishing the bevel faces in 2 steps and to present a lens end face machining method, wheel and device for eyeglass lens end face machining with which it is possible to speed up the polishing process and make finishing precision uniform and obtain fashionable eyeglass.\nMoreover, in addition to the aesthetic problem of the lens end face remaining white after bevel polishing that was previously described, there is a problem with polishing precision and fashionable eyeglasses in that when planing, such as smooth machining and machining to a mirror polishing etc., is performed with a wheel that has a bevel-groove and a planing face, streaks are made. That is, cylindrical grinding stone called diamond wheels(stone) have a bevel-groove for formation of a bevel at the end face of the eyeglass lens and a flat face for flat machining the end face of an eyeglass lens. In further detail, the wheel has groove inclined face 301 for V finishing having a specific angle with respect to the axial direction called angle No. 1, flank 203 for the eyebrow of the frames continuous with this groove inclined face 301 having a specific angle with respect to the axial direction referred to as angle No. 2 that is smaller than angle No. 1, and flat finishing face 303 continuous with this flank 302 for flat machining and parallel to the axial direction on the surface around the periphery of the wheel. The inclination at boundary K between above-mentioned flank 302 and flat finishing face 303 is not continuous.\nConsequently, when an eyeglass lens moves past boundary K to the left in the direction of the X axis during flat machining, apex A of the end face of eyeglass lens 6 straddles boundary K and a streak from boundary K is made in end face 6a of eyeglass lens 6. When a streak is made in end face 6a of the eyeglass lens, edging precision drops and becomes non-uniform, and the product is not fashionable. Therefore, such a streak is undesirable. This is particularly a problem with flat finished surfaces that remain white and are further given a mirror finish so that they are transparent.\nThereupon, in order to solve this problem, flat finishing face 303 is made longer in the axial direction so that even if eyeglass lens 6 moves to the left in the direction of the X axis during flat edging, it will not pass boundary K. However, there is a problem in that as a result, wheel 1 is larger.\nInerdentally there is a demand for mechanical polishing of the bevel face that remains white using a wheel as a means of solving the above-mentioned aesthetic problem of the lens end face remaining white after bevel polishing because buffing, etc., manually takes time and increases cost. However, there is also a problem when a polishing wheel is used with the existing wheel in that the device becomes bigger.\nThe objective of the present invention is to solve the above-mentioned problems of prior art and present an eyeglass lens end face matching method with which polishing precision is uniform, the product is excellent in terms of being fashionable, and the device can be reduced in size. Another objective of the present invention is to present an eyeglass end face machining method with which it is possible to add a polishing wheel that can give the eyeglass lens end face a mirror polish without greatly increasing length of the wheel in the axial direction."}
{"text": "The present invention generally relates to an image recording apparatus with a scanner unit, and more particularly to an image recording apparatus with a scanner unit which selectively performs a printer mode in which an image is printed on a copy sheet, and a scanner mode in which an image is read from a document. This apparatus is applicable to laser printers, copying machines, facsimile machines and others.\nConventionally, a printing apparatus with a scanner unit can perform a scanner mode and a printer mode. When the scanner mode is performed by this printing apparatus, an original document is passed through a scanner unit and an image is read from the document by illuminating the document with light emitted by a light source of the scanner unit so that a signal indicating the original image is produced from the light reflected from the document.\nIn the meantime, when the printer mode is performed by the printing apparatus mentioned above, a copy sheet is fed from a paper cassette in the printing apparatus, and the copy sheet is transported to a photosensitive drum via a pair of registration rollers. The photosensitive drum is electrostatically charged by a charging unit, and the surface of the photosensitive drum is exposed to light emitted by a light source of an optical recorder unit, so that an electrostatic latent image in accordance with an image signal is formed on the photosensitive drum. By means of the pair of registration rollers, the time of forming the electrostatic latent image on the photosensitive drum is harmonized with the time of transporting the copy sheet to the photosensitive drum.\nFIG. 1 shows the construction of the optical recorder unit provided within the above mentioned printing apparatus. As described above, at this optical recorder unit, an electrostatic latent image is formed on a photosensitive medium. A light beam emitted by a laser diode 1 of the optical recorder unit in accordance with an image signal modulated by a driver 7 is passed through a collimator lens 2 and a cylindrical lens 3, and it is deflected by a polygonal mirror 4 which is rotated by a driving motor (not shown) around a rotating axis of the polygonal mirror 4. The deflected light beam is passed through a focusing lens 5 and a troidal lens 6 to convert the light beam into a converging light beam. This converging light beam is cast on a photosensitive drum 8 to form a beam spot on the photosensitive drum 8. Because the polygonal mirror 4 is rotated by the driving motor around the rotating axis of the polygonal mirror 4, the photosensitive drum 8 is scanned by the converging light beam from the polygonal mirror 4 in a main scanning direction indicated by an arrow A in FIG. 1. The photosensitive drum 8 is also rotated around its rotating axis, and it is scanned by the converging light beam in a sub scanning direction perpendicular to the main scanning direction. Thus, an electrostatic latent image is formed on the photosensitive drum 8 by the optical recorder unit in accordance with the image signal.\nIn the printing apparatus described above, a developing unit including a toner cartridge is provided. When the electrostatic latent image on the photosensitive drum is passed through the developing unit, the electrostatic latent image is developed on the photosensitive drum with toner, so that the electrostatic latent image is made visible. This toner image is transferred by a transfer charging unit from the photosensitive drum to a copy sheet. This copy sheet has been transported to the photosensitive drum 8 via the registration rollers in a synchronized manner. Then, the copy sheet with the toner image is transported from the photosensitive drum to a fixing unit including a fixing heater, so that the toner image is fixed to the copy sheet by the heat and pressure of the fixing unit. The copy sheet is then transported from the fixing unit to a copy sheet ejection outlet part of the printing apparatus. Each copy sheet transported from the fixing unit is stacked at the copy sheet ejection outlet part.\nHowever, in the above described printing apparatus, during the scanner mode operation in which the scanner unit is operated, some print processing units including the optical recorder unit, the fixing unit and the transfer charging unit, which should be used in the printer mode only, are unnecessarily driven. Thus, in the above described printing apparatus, there is a problem in that the driving of such printing processing units for a long period of time during the scanner mode operation will shorten an operational life of each of the units in the printing apparatus, and also the electric energy consumed by the printing apparatus is higher than a reasonable energy consumption level. Also, the construction of the conventional printing apparatus is bulky because both the driving mechanism of the built-in scanner and the driving mechanism of the print processing units are necessary.\nMoreover, at the time of an original document sheet being passed through the transfer charging unit or the fixing unit when the scanner mode is performed in the printing apparatus described above, there is a problem in that an original document to be scanned may be damaged due to the pressure of the transfer charging unit or the heat and pressure of the fixing unit, and that a residual toner on the photosensitive medium is transferred to the original document sheet."}
{"text": "In mass transfer processes such as waste treatment and bio-reactions, it is common to carry out these processes in an aeration vessel in which gas, such as oxygen and/or air, is introduced into a biodegradable liquid for treatment. These aforementioned processes are oftentimes utilized by municipalities and industry to treat waste water wherein the object of the process is to introduce air to the liquid and then micro-organisms in the liquid proceed to use this oxygen to digest the waste. The gas is commonly introduced by way of impellers wherein the impellers aerate the liquid.\nDuring the aeration of a liquid, for example, waste water treatment, it is common to employ impellers which are especially adapted for use in surface aeration of liquids in an open tank where the impellers are disposed at the surface of the liquid in vessel. Typical surface aerators commonly used in the art are generally either radial flow impellers or pitched blade turbines. The blades are usually flat rectangular plates which are pitched usually at an angle of 45° to the axis of rotation of the impeller. This is frequently termed the static level of the liquid. The aforementioned impellers are commonly located close to the static level liquid surface and a small portion of the width of the blade may project up through the surface. Typically, when the impeller is pitched forwardly, the upper edge of the blade is termed the leading edge while lower edge is termed the trailing edge. Alternatively, typically when the impeller is pitched backwardly, the upper edge is the trailing edge while the lower edge is the leading edge. The liquid is usually either pushed out in front of the angled blade and/or scooped by the blade and discharged radially across the surface of the tank with some of the liquid being sprayed into the atmospheric air from the outer upper surfaces of the blade. As a result of the spraying of the liquid into the atmospheric air, the liquid becomes aerated."}
{"text": "With the proliferation of the automated interactive machines, exemplified by the automated teller machines (ATM) for financial transactions, there has been an emerging need for a more reliable personal identification system for authenticating users who desire to conduct transactions remotely and automatically without human intervention. Conventionally, a person simply inserts her ATM card into the machine to have her account information and password, or PIN (\"Personal Identification Number\", used here interchangeably with the word \"Password\"), read. However, as the everyday life as a whole becomes more automated and security-conscious, a person often has to manage various different passwords and PIN's, for accesses to her banking account, her home security system, or her eMail account, to name just a few. This overflow of information has already contributed to the complexity of conventional personal identification systems in that without the correct password for an ATM, a legitimate user may be denied of her access to her account or her on-line brokerage account.\nThere is an often overlooked burden placed on the institutions providing on-line, or remote, transactions which are accessed through the customers' passwords or PINs. Maintaining the passwords or PINs forces the financial institutions to allocate additional machines and human resources to manage interface with customers when a customer forgets her Pin or when a customer requests her PIN be changed.\nAlso, passwords have been proven to be insufficient in preventing fraud, where all a would-be criminal needs is an ATM card and the password, which are both reasonably within the reach of those unscrupulous ones. This is just the first example of how the conventional personal identification paradigm is vulnerable, in addition to being complex as discussed above.\nAnother problem plagues the integrity of the supposedly secured financial transaction, where sometimes it is the actual account holder who defrauds the institution by first accessing her account and later denying such transaction from ever taking place. While there is a limit as to the extent of this sort of heinous behavior, it amounts to a significant sum even with just a small percentage of the ATM transactions considered. Without a more reliable identification system, institutions will just have to write off the losses or pass the losses to the rest of the consumers, thereby increasing everyone's cost of doing business.\nAside from the ATM transactions, with the increasing affordability, as well as sophistication, of personal computers and telecommunication hardware and software, it is more likely that one will soon be accessing a host of information or conducting a variety Of secured transactions using a PC, a modem and a common public switching network, such as Prodigy and Internet, etc. Authentication thus becomes an even more paramount task for the industry to tackle.\nA simple personal identification system may address the above problems. Fingerprints have been known years ago to have a high degree of accuracy and reliability. One never forgets her fingerprints, or confuses the fingerprints with other information. Also, a criminal cannot steal or duplicate someone's fingerprints to impersonate the account holder, generally speaking. Therefore, fingerprints are essentially a personal identification with a one-to-one correspondence, given that the fingerprint recognition systems have progressed along with the information revolution. Companies such as Identix and Startech have developed front-end fingerprint image recognition systems to reliably and accurately analyze and recognize fingerprints.\nAt the back-end, major processor suppliers such as IBM and AT&T already have systems in place to provide a linkup with the fingerprint image recognition systems such that the massive fingerprint database may be linked and accessed for the institution to quickly authenticate the person in front of its machine, or the person seeking to access her brokerage account through a PC with a modem. To a certain extent, the present front-end and back-end suppliers have reached a point, where it is merely a matter of time before their capabilities and achievements can be fully utilized by the industry, especially the financial industry.\nEven with reliable fingerprint image recognition systems at the front-end and quick-response processor at the back-end, there are still problems with this paradigm. Assuming it is reasonably affordable for a PC owner to have a personal fingerprint recognition device to provide access to her on-line brokerage account at a brokerage firm with a processor to facilitate authentication, there is still about 1% error rate, generally characterized by false rejection of legitimate users, due to the inherent imperfection of one's fingerprints. For example, if a person regularly works with abrasive chemicals, the quality of her fingerprints tends to deteriorate throughout the years. The degraded quality of the fingerprints, when faced with a security sensitive system as in most security-sensitive transactions, will certainly add to the agony of the users, thus further eroding the public's confidence toward the integrity of future systems.\nOn the other hand, if the security sensitivity is forced to be compromised to minimize false \"rejection\", then the error rate of false \"acceptance\" may increase and vice versa. Conversely, if the security sensitivity is forced to be compromised to minimize false \"acceptance,\" then the error rate of false \"rejection\" may increase. Now that a half-way decent \"match\" will allow access erroneously. This is also not something which will contribute to the public's confidence toward fingerprint-based personal identification systems. Nor will it contribute to the industry whose primary application of the fingerprint-based personal identification systems is to protect their business and financial interests.\nFurthermore, the creation of an initial file, i.e., when the account holder first sets up her account with her fingerprints at the institution's facility, may not be perfectly analyzed and stored as file data. The possibility of having less than perfect fingerprints on file makes the occurrence of false rejection/acceptance even more likely. For example, if the initial registration has a 90% accuracy, it would always be a 90% accuracy. It would still be a 90% match at best, even with a 100% accurate reading at the ATM at a later time. In other words, both ends of the overall system may contribute to the unreliability of the system.\nTherefore, it is desirable to have a personal identification system for use with fingerprint recognition front-ends to raise the percentage of accuracy, thus minimizing the security risks in connection with secured transactions.\nIt is also desirable to have a personal identification system for taking advantages of the conventional fingerprint recognition devices to provide a flexible solution in light of the various vendors of the front-end and back-end systems.\nIt is further desirable to have a fingerprint-based personal identification system which will provide an easy-to-use solution to the security issues involved in accessing the information superhighway."}
{"text": "Steam turbines include a shell that functions to contain the high pressure, high temperature steam and to support the nozzles and casings that direct steam in the most efficient manner possible through rotating airfoils to produce maximum torque on the shaft. The shell includes support arms that extend outward from the shell. Shell arms rest on a support structure that is integral to a structure that also provides support for the turbine rotor, and which serves to house other turbine related components and instrumentation. This structure is often referred to as a “standard”. The elements of the stator portion of the steam turbine are coupled to the shell, thus, the stator portion of the steam turbine is supported by the shell arm support structure.\nThe rotor of the steam turbine is typically supported by bearings. The bearings are typically mounted within a bearing housing that is supported by a bearing support structure. The bearing support structure can be a part of the “standard” mentioned above. The bearing support structure, while it may be integral with the standard, is not integral with the shell arm support structure. The shell arm support structures and the bearing support structures are exposed to different environmental conditions during various stages of operation of the steam turbine. During transient operations, including but not necessarily limited to startup, load changes, shutdown, and cool downs while on turning gear, different portions of the steam turbine and the supporting elements experience changes in temperature. These changes in temperature may occur at different rates within the different parts of the steam turbine and its support structure, which leads to differential thermal growth of the steam turbine elements, and the support elements.\nDuring startup operations, the bearings and the bearing support structure which supports the bearings and the rotor of the steam turbine tend to increase in temperature more quickly than the shell arm support structure. In part, this occurs because the bearings supporting the rotor rapidly heat up during startup because they are being driven by the changing oil temperature which they are in constant contact with, and the heat generated in the bearings is transferred to the bearing support structure. In contrast, the shell and the shell arm support structure, which are not in constant contact with the oil, tend to warm up more slowly.\nSimilarly, the bearings typically cool more rapidly upon shutdown because oil supply temperature is lowered as the steam turbine moves from full speed operation to operation at turning gear speeds. In contrast, the support structure beneath the shell arms generally cools very slowly because its temperature is driven more by shell temperature and conduction of that heat from the shell arms into the structure. Shell temperatures decay very slowly. And this can cause the bearing support structure to cool more quickly than the shell arm support structure upon shutdown.\nWhen there are differences in the rate at which the temperature of the bearing support structure increases or decreases, as compared to the rate at which the temperature of the shell arm support structure increases or decreases, the temperature differences can lead to differences in the rate or amount of thermal expansion and contraction of these two elements. And differences in the rate or amount of thermal expansion or contraction as between the bearing support structure and the shell arm support structure can cause changes in the amount of radial clearance available between rotating and stationary parts."}
{"text": "As computers have grown increasingly important in today's society, the importance of public and private networks and, especially, the Internet has also increased. As increasing numbers of users access the Internet, the need for efficient use of bandwidth has also increased. The increasing numbers of requests handled by the Internet are increasing the delay experienced by a user between generating a request and receiving a response to the request because of bandwidth limitations.\nOne traditional solution to decreasing overall bandwidth usage and decreasing the delay experienced by the user has involved caching previously requested content at the user's computer for faster retrieval. A related traditional solution has involved caching previously requested content for multiple users at a single cache server. Another traditional solution has involved increasing the bandwidth of the network connection between the Internet, the user and the web servers handling the requests. However, traditional solutions have often failed as the number of requests continue to increase and overload single cache servers and because of the expense associated with maintaining large numbers of high speed connections to the Internet. In addition, traditional solutions have often failed to provide for the distinguishing the relative importance of web pages."}
{"text": "1. Field of the Invention\nThe present invention relates to a storage device, a storage/reproducing device, and a rewritable reproducing device that are incorporated in or attached to a personal computer as a unitary structure. More particularly, the invention relates to a storage device, a storage/reproducing device, and a reproducing device that are incorporated in or attached to a desk-top personal computer as a unitary structure.\n2. Background of the Invention\nThe storage device of this kind has heretofore not been equipped with an external data input means, and data have been always input through a personal computer such as a desk-top personal computer one. Without a personal computer, therefore, the storage device is unable to operate. However, with the widespread use of various information equipment such as PHS, digital camera and the like, which have a prerequisite of transferring the data to a personal computer, it would be more convenient if transfer of data can be relayed by a storage device.\nTherefore, there arises a technical problem to be solved that storage devices of this kind must be provided with a function of directly fetching the data from an external unit and relaying the transfer of data. The object of the present invention is to solve this problem.\nSo far, furthermore, conventional storage/reproducing devices and the reproducing devices of this type have usually been formed in the shape of a thin box. Therefore, a storage/reproducing device of built-in type has a narrow area exposed from the front of the desk-top personal computer, and is difficult to secure a space where the connector for fetching record data from an external data unit is installed. It may be possible to lead out the connector cable from the back of the personal computer and lay the cable along the outer side of the personal computer. However, this causes the operation to be complex, marring the appearance. It would be convenient if the storage/reproducing device and the read-only reproducing device of the attached type, too, have a function of easily fetching data from an external data unit.\nTherefore, there arises a technical problem to be solved that storage/reproducing devices and the reproducing devices of this kind are improved to easily fetch external data. The object of the present invention is to solve this problem."}
{"text": "A large number of analytical operations benefit from the simultaneous illumination of relatively large area of substrates in order to accomplish the desired analysis. For example, interrogation of biopolymer array substrates typically employs wide area illumination, e.g., in a linearized beam, flood or reciprocating spot operation. Such illumination allows interrogation of larger numbers of analytical features, e.g., molecule groups, in order to analyze the interaction of such molecule groups with a sample applied to the array.\nIn the case of DNA arrays, in particular, large numbers of oligonucleotide probes are provided in discrete locations on a planar substrate surface, such that the surface comprises multiple, small patches of identical probes, where the probes' nucleotide sequence for each patch location is known. When one applies a labeled sample sequence to the array, the position on the array to which the sample hybridizes is indicative of the complementary probe sequence, and as such, the sequence of the sample sequence. These arrays are generally interrogated using laser based fluorescence microscopes that are capable of applying excitation illumination over large areas of the substrate in order to interrogate all of the patches. Such systems have employed galvo scanners, slower, scanning microscopes, linearized beam illumination, and wide area flood illumination.\nIn some cases, however, a more tightly controlled illumination strategy may be desired. For example, it may be desirable to provide stricter control of the volume of material that is illuminated, as well as the overall area that is illuminated, effectively controlling illumination not only in one of the x or y axes of a planar substrate, but also in the z axis, extending away from the substrate. One example of controlled illumination that accomplishes both lateral (x and y) and volume (z) control is the use of zero mode waveguides as a base substrate for analyzing materials. See, U.S. Pat. Nos. 6,991,726 and 7,013,054, the full disclosures of which are hereby incorporated herein by reference in their entirety for all purposes. Briefly, zero mode waveguide array substrates employ an opaque cladding layer, e.g., aluminum, chromium, or the like, deposited over a transparent substrate layer, and through which are disposed a series of apertures through to the transparent layer. Because the apertures are of sufficiently small cross sectional dimensions, e.g., on the order of 50-200 nm in cross section, they prevent propagation of light through them that is below a cut-off frequency. While some light will enter the aperture or core, its intensity decays exponentially as a function of the distance from the aperture's opening. As a result, a very small volume of the core is actually illuminated with a relevant level of light. Such ZMW arrays have been illuminated using a number of the methods described herein, including spot illumination, flood illumination and line illumination (using a linearized beam)(See, e.g., co-pending U.S. patent application Ser. No. 11/483,413 (filed Jul. 5, 2006), and 60/772,908 (filed Feb. 13, 2006), the full disclosures of which are incorporated herein by reference in their entirety for all purposes).\nWhile the various foregoing systems and methods have proven some measure of effectiveness, the present invention provides for improvements over these systems and methods, in a number of respects."}
{"text": "The present invention relates to a method for adjusting pressure in a fuel tank and a corresponding fuel tank system.\nThe following discussion of related art is provided to assist the reader in understanding the advantages of the invention, and is not to be construed as an admission that this related art is prior art to this invention.\nIt is desirable to reduce the release of hydrocarbons from automobiles with an internal combustion engine into the environment. To this end, automobiles typically have an active charcoal filter which is generally arranged between a tank of the automobile filled with combustible fuel and the environment. This active charcoal filter operates as temporary storage for evaporated hydrocarbons from the automobile. The active charcoal filter is typically regenerated in regular intervals with a flow of fresh air, which takes up the hydrocarbons stored in the active charcoal filter and transports them to the internal combustion engine. These so-called purge phases of the active charcoal filter are intended to ensure that the filter has always sufficient storage capacity for storage of hydrocarbons.\nThese purge phases cannot be performed very frequently particularly in hybrid vehicles, because the internal combustion engine is only operated intermittently. To ensure that the active charcoal filter has always enough capacity for absorbing hydrocarbons, it has been proposed to reduce the quantity of hydrocarbon vapor generated in the tank by adjusting a predetermined internal tank pressure greater than ambient pressure. To this end, an electrically controllable shutoff valve is arranged between the tank and a connecting line which connects the tank with the active charcoal filter. When the tank is heated, the pressure inside the tank increases to a predetermined, presettable threshold value for the pressure. Hydrocarbons in the fuel then increasingly remain in their liquid phase. The shutoff valve is opened only when this threshold value is exceeded, whereafter hydrocarbon-containing vapors escape from the tank. These vapors are then transferred through the connecting line to the active charcoal filter, where they are absorbed, causing the pressure in the tank to decrease again. As soon as the pressure in the tank is again below the threshold value, the shutoff valve is closed again.\nDisadvantageously, the pressure inside the tank cannot be adjusted to be variable. A certain threshold value for the pressure must be preset in advance. A very high pressure is typically selected as the threshold value, allowing vapors to escape from the tank as infrequently as possible. The shutoff valve is used only as an emergency valve in order to prevent damage caused by a very high pressure. The complete tank system therefore requires a rather complex structure to ensure that the tank can withstand regularly occurring high pressures. The manufacturing costs of such tank systems are therefore very high. The absorption capacity of the active charcoal filter is also sometimes not fully utilized.\nIt is therefore an object of the present invention to provide a method for adjusting pressure in a fuel tank and a fuel tank system, with which the pressure in the fuel tank can be variably adjusted and with which the manufacturing cost of the tank system can be reduced.\nIt would therefore be desirable and advantageous to obviate prior art shortcomings and to provide an improved method for adjusting pressure in a fuel tank and a fuel tank system with which the pressure in the fuel tank can be variably adjusted, resulting in a reduction of the manufacturing costs of the tank system."}
{"text": "Maneuvering spacecraft in the vicinity of other objects is a complex operation, especially when the objective is to rendezvous with another, target object in a docking mission. Most of the algorithms developed to control a space craft in such circumstances have serious drawbacks.\nMost of the existing methods and systems do not, for instance, optimize the fuel cost of the maneuver, or minimize the time taken to perform the maneuver. Moreover, they cannot, in general, deal with having unanticipated additional objects in the vicinity. Nor can they deal with stayout zones that have to be avoided that may result from, for instance, the optical sensors on the target. Most of the existing methods cannot be used to orbit the target en-route to a rendezvous, which may be desired in order to inspect either the spacecraft or the target. The few algorithms that can deal with obstacles and stay out zones tend to be very computationally intensive, requiring offline calculation. They are, therefore, of limited use, especially if obstacles near the target change during the maneuver. This may happen, for example, if the target encounters unexpected additional targets.\nSeveral methods have been developed for use by spacecraft for close maneuvering, rendezvous and docking missions. Four of these are discussed below.\nMethod 1. The glideslope algorithm is a popular method of approaching another spacecraft. In this method the range rate is kept proportional to the range. When the range is zero the range rate is zero. The initial range rate/range ratio can be any desired. This algorithm cannot deal with stayout zones or obstacles, thus these must be handled in an ad hoc fashion. In addition, it does not allow for orbiting the target.\nMethod 2. Another method is the rbar or vbar approach. Vbar is the axis along the velocity vector and rbar is the axis along the position vector to the Earth. Approaches are made along either of these axes. However, this method also does not account for stayout zones or obstacles. As with the glideslope algorithm, this method does not handle orbiting the target.\nMethod 3. A method developed by Miller that can handle trajectory constraints such as obstacles and stayout zones is parametric programming and model predictive control, which finds the optimal control inputs as a piecewise affine function of the states. Most computation is performed off-line thus reducing the on-line computation to a simple table lookup. The solution found off-line is an exact solution, the equivalent of solving an open-loop optimal control problem over a finite horizon at each time step. However, even though the computation is offline it still must be done on the flight processor. In addition, obstacles may change during the approach which means this computationally intensive procedure must in practice be done in real-time.\nMethod 4. Another method developed by Zhao and Yang employs a 4-dimensional search space with 3 spatial and 1 temporal coordinate. Obstacles and potential conflicts are represented by several basic shapes and linear combinations of these shapes. Mathematical conditions are developed for a given point as well as a trajectory segment between two points to be outside of an obstacle. The A* search technique is used to obtain trajectory solutions in which successor trajectory points are selected that both avoid obstacles and satisfy dynamic motion constraints of the vehicle. A linear combination of flight distance and time is optimized. This algorithm is able to generate flight trajectories rapidly and has the potential to be used in real-time. However, the used of simple geometric shapes as the basis for obstacles is a major limitation of this approach. In addition it has no provision for non-uniform node spacing and generation."}
{"text": "1. Field of the Invention\nThe present invention relates generally to text messaging communications, and more particularly, to graphical expression during text messaging communications.\n2. Description of the Prior Art\nText messaging refers to the exchange of brief written messages over a network between phones and/or computing devices. Mobile instant messaging (MIM) technology extends text messaging service accessibility to mobile computing platforms (e.g., standard mobile phones, smartphones, and electronic tablets). With MIM technology, instant messaging services can be accessed from computing devices, including standard mobile phones and smartphones using a myriad of operating systems.\nIn text messaging, whether fixed-line or mobile, real-time text messages are communicated directly between individuals using computing devices (e.g., personal computers or mobile phones). Two types of text messaging are instant messaging and online chat. Although qualitatively similar, instant messaging is used in common parlance to refer to communications between known users (e.g., using a contact list or a friend list) whereas online chat is used to refer to communications among anonymous users in a multi-user environment.\nBoth instant messaging and online chat use a push-based style of internet communication. As with email, a user types in text that is transferred over a network (e.g., the Internet) to another user. Text messages between users are perceived to be quasi-synchronous—that is, to occur nearly simultaneously. This near-simultaneity is due to a push-based style of interne communication (i.e., communications are “pushed” to recipients by a central server when received) which allows immediate acknowledgement of, or reply to a message and makes text messaging an efficient and effective means of communication between users.\nOne inherent limitation of text messages is that each message is limited to 160 bytes of data (although some providers reserve some of those bytes for service use), or approximately 160 alphanumeric characters of the English alphabet. This size limitation severely restricts the type of data that can be transmitted via text message. As a result, graphical expression in text messages has been limited to one or more textual emoticon composed of alphanumeric symbols (e.g., a colon and a right parenthesis to symbolize a happy face, or a smiley face included within a font set), an attached graphic image file (e.g., a .GIF file of a winking smiley face or a flash animation of kissing lips), or a simple descriptive statement of an imaginary graphic action (“John has thrown a cow at you.”). Because these graphical options are limited, the ability to effectively express emotion within text messages is minimal, at best.\nPrior art messaging systems do allow the user to upload a personal picture and/or a tag phrase (or sentence) to his profile which is/are seen by other users during messaging sessions. The personal picture and tag phrase offer the user an opportunity to express a snapshot of his thoughts and/or mental state—at least from the instant in time that the profile was created—to others with whom he is communicating. Current messaging systems, at most, have the capacity to allow the user to change his associated profile picture during the text messaging session. Although the user can change his profile picture to express a different state or impression, the change is implemented by uploading a new picture outside of the text messaging session."}
{"text": "In general, a hybrid electric vehicle (HEV) is a vehicle that uses two types of power sources together, and the two types of power sources are typically an engine and an electric motor. Such a hybrid vehicle has excellent fuel efficiency and power performance and is advantageous in that the amount of exhaust emissions is reduced compared to a vehicle having only an internal combustion engine, and thus has been actively developed in recent years.\nFIG. 1 illustrates an exemplary structure of a powertrain of a general hybrid vehicle.\nReferring to FIG. 1, the powertrain of the hybrid vehicle adopts a parallel-type hybrid system, in which an electric motor (or a drive motor) 140 and an engine clutch (EC) 130 are mounted between an internal combustion engine (ICE) 110 and a transmission 150. In particular, since the electric motor 140 is mounted close to the transmission 150, the powertrain system of the hybrid vehicle may also be referred to as a transmission-mounted-electric-device-type (TMED-type) system.\nTypically, when a driver depresses an accelerator pedal after starting the vehicle, the motor 140 is first driven using the electrical power of a high-voltage battery 160 in the state in which the engine clutch 130 is opened, and wheels are moved by the power transferred from the motor to a final drive (FD) (not shown) via the transmission 150 (i.e. an EV mode). When a greater driving force is required due to the gradual acceleration of the vehicle, the engine 110 may be driven by operating an auxiliary motor (or a starter/generator motor) 120.\nThus, when the number of revolutions per minute of the engine 110 and the number of revolutions per minute of the motor 140 are equal to each other, the engine clutch 130 is engaged so that the vehicle is driven by both the engine 110 and the motor 140 or by only the engine 110 (i.e. transition from the EV mode to an HEV mode). When a predetermined engine-off condition, such as the deceleration of the vehicle, is satisfied, the engine clutch 130 is opened and the engine 110 is stopped (i.e. transition from the HEV mode to the EV mode). In such a hybrid vehicle, a battery may be charged by converting the driving force of the wheels into electrical energy when a braking operation is performed, which is referred to as braking energy regeneration or regenerative braking.\nThe starter/generator motor 120 serves as a starter motor when the engine is started, and also serves as a generator when the rotational energy of the engine is recovered after starting or when starting off. Therefore, the starter/generator motor 120 may be referred to as a “hybrid starter generator (HSG)”, or may also be referred to as an “auxiliary motor” in some cases.\nDescribing the HSG 120 and the electric motor 140 in terms of current flow, the high-voltage battery 160 outputs a DC voltage, and an inverter 180 appropriately converts the DC voltage into an AC voltage in accordance with a torque command and a rotational speed of the electric motor 140. The converted AC voltage is supplied to the electric motor 140 and the HSG 120 so as to drive the vehicle or start the engine. The inverter 180 also converts the counter electromotive force of the electric motor 140 and the HSG 120 due to the driving force of the traveling vehicle or regenerative braking into a DC voltage so as to charge the high-voltage battery 160.\nHereinafter, the flow of high-voltage energy depending on whether the electric motor is performing a charge or discharge operation will be described with reference to FIG. 2.\nFIG. 2 is a view schematically illustrating the flow of energy depending on the operation of the electric motor in a general hybrid vehicle.\nReferring to FIG. 2, in the hybrid vehicle, when the electric motor 140 operates in a discharge mode, electrical energy for driving is transferred from the battery 160 to the electrical motor 140 via the inverter 180. When the electric motor 140 operates in a charge mode, charge energy, which is generated by the electrical motor 140, is transferred to the battery 160 via the inverter 180.\nOn the assumption that the electric motor 140 rotates at a constant speed in the positive (+) direction, the motor torque has a positive (+) value in the discharge mode for generating driving force, and has a negative (−) value in the charge mode.\nHowever, at some operating points of the motor, the loss of energy may be larger than the charge energy generated by electricity generation due to loss in the energy transfer path, including copper loss of the motor. A driver expects that the battery will be charged whenever regenerative braking is performed, but in practice there is a problem in that the battery may be discharged depending on the operating point of the motor."}
{"text": "1. Technical Field\nVarious embodiments of the present invention relate to a data storage device and an operating method thereof.\n2. Related Art\nRecently, the computing environment paradigm has converted to ubiquitous computing so that computer systems can be used anytime and anywhere. Due to this fact, the use of portable electronic devices such as mobile phones, digital cameras, and notebook computers has rapidly increased. In general, such portable electronic devices use a data storage device having a memory device. The data storage device is used as a main memory device or an auxiliary memory device of the portable electronic devices.\nA data storage device using a memory device provides advantages of excellent stability and durability, fast information access, and low power consumption. Data storage devices having such advantages include universal serial bus (USB) memory devices, memory cards having various interfaces, universal flash storage (UFS) devices, and solid-state drives (SSD).\nAs large capacity files such as audio files and video files are used in portable electronic devices, data storage devices need to have large storage capacities. Therefore, the data storage device includes a plurality of memory devices. To efficiently control or access the memory devices, the data storage device may set memory regions of the respective memory devices as access units."}
{"text": "The use of payment cards, such as credit cards, debit cards, and pre-paid cards, has become ubiquitous. Most payment card accounts have one or more associated physical cards; however, the use of non-traditional payment devices, such as appropriately-configured “smart” cellular telephones, is increasing.\nTax-advantaged financial accounts of various kinds are quite prevalent. One non-limiting example includes a flexible spending account (FSA), also known as a flexible spending arrangement, which allows an employee to set aside a portion of earnings, pre-tax, to pay for qualified expenses such as medical, dependent care, or the like. Another non-limiting example is a health savings account (HSA), which is a tax-advantaged medical savings account available to taxpayers in the United States who are enrolled in a high-deductible health plan (HDHP)."}
{"text": "1. Field of the Invention\nThis invention relates generally to toys and accessories therefor, and more specifically to an apparatus for filling toy balloons with a liquid such as water.\n2. Description of the Prior Art\nToy balloons are typically made from latex rubber or similar expandable material, and are usually filled with air or other gas for use as a toy or decoration. Alternatively, many users desire to fill their balloons with water or other liquid so they may use them as \"water balloons\", which can then be thrown to intentionally break or \"explode\" upon impact with a target as such as the ground, or an unfortunate playmate. It is the process of filling these water balloons with water that is the subject of this invention.\nMost users fill a water balloon by expanding the balloon opening (neck) over a sink faucet, hose bib, drinking fountain, or other pressurized water supply fitting. The user then opens the water supply valve so that the pressurized water fills the balloon with water to the desired volume, whereupon the user closes the water supply valve, removes the filled balloon from the fitting, and ties off and seals the neck of the water balloon for play.\nHowever, many play areas do not have access to a supply of water under pressure, making it difficult if not impossible to fill water balloons. For example, while swimming pools, creeks, lakes, oceans, and other bodies of water make ideal locations for water balloon play, it is difficult to fill the balloons with water for lack of a pressurized water supply."}
{"text": "Thermal imaging or thermography is a recording process wherein images are generated by the use of imagewise modulated thermal energy.\nIn thermography two approaches are known:\n1. Direct thermal formation of a visible image pattern by imagewise heating of a recording material containing matter that by chemical or physical process changes colour or optical density.\n2. Formation of a visible image pattern by transfer of a coloured species from an imagewise heated donor element onto a receptor element.\nA survey of \"direct thermal\" imaging methods is given in the book \"Imaging Systems\" by Kurt I. Jacobson-Ralph E. Jacobson, The Focal Press--London and New York (1976), Chapter VII under the heading \"7.1 Thermography\". Thermography is concerned with materials which are not photosensitive, but are heat sensitive. Imagewise applied heat is sufficient to bring about a visible change in a thermosensitive imaging material.\nAccording to a direct thermal embodiment operating by physical change, a recording material is used which contains a coloured support or support coated with a coloured layer which itself is overcoated with an opaque white light reflecting layer that can fuse to a clear, transparent state wherein the coloured support is no longer masked. Physical thermographic systems operating with such kind of recording material are described on pages 136 and 137 of the above mentioned book of Kurt I. Jacobson et al.\nYet most of the \"direct\" thermographic recording materials are of the chemical type. On heating to a certain conversion temperature, an irreversible chemical reaction takes place and a coloured image is produced.\nIt has been suggested to use a thermoreducable silver source in combination with a reducing agent in a direct thermal film in order to increase the optical density in transmission of a printed image (see EP-A-537.795). Although continuous tones can be obtained by said printing method, the gradation produced by said printing method is too high resulting in only a few intermediate density levels. Fluctuations in the heat transfer from the heat source to the printing material result in a density difference of the final image. Thus, it is extremely difficult to obtain images having a uniform density profile. A direct thermal printing method moreover has the disadvantage that in the non-image places the co-reactants always remains unchanged, impairing the shelf-life and preservability.\nThermal dye transfer printing is a recording method wherein a dye-donor element is used that is provided with a dye layer wherefrom dyed portions or incorporated dye is transferred onto a contacting receiving element by the application of heat in a pattern normally controlled by electronic information signals.\nIn European Patent Application No. 94200612.3, a thermal imaging process is provided using\n(i) a donor element comprising on a support a donor layer containing a binder and a thermotransferable reducing agent capable of reducing a silver source to metallic silver upon heating and (ii) a receiving element comprising on a support a receiving layer comprising a silver source capable of being reduced by means of heat in the presence of a reducing agent, said thermal imaging process comprising the steps of\nbringing said donor layer of said donor element into face to face relationship with said receiving layer of said receiving element, PA1 image-wise heating a thus obtained assemblage by means of a thermal head, thereby causing image-wise transfer of an amount of said thermotransferable reducing agent to said receiving element in accordance with the amount of heat supplied by said thermal head and PA1 separating said donor element from said receiving element. PA1 bringing said donor layer of said donor element into face to face relationship with said receiving layer of said receiving element, PA1 image-wise heating a thus obtained assemblage preferably by means of a thermal head, thereby causing image-wise transfer of an amount of said thermotransferable reducing agents to said receiving element in accordance with the amount of heat supplied and PA1 separating said donor element from said receiving element.\nThis printing method is further referred to as `reducing agent transfer printing` or `RTP`.\nHowever, the stability of the donor element in said European Patent Application has been found to be poor. More particularly, the reducing agent tends to crystallize in the donor layer. As a result of this crystallization during storage, transfer of reducing agent is seen during printing, even on places where no heat has been applied by the thermal head. This leads to a printing fog in the final image. This problem is especially seen when a high amount of reducing agent is used in the donor layer. This high concentration is necessary to obtain high optical densities of the final printed image (above 2.0-2.5).\nMoreover, the neutral hue of the grey scale of the printed image is dependent on the choice of a specific reducing agent. It is extremely difficult to find reducing agents that yield a neutral grey-tone (e.g. for medical diagnostics). It has also been found that the reducing agent shows the disadvantage after transfer that the oxidation product of the reducing agent tends to crystallize in the receiving element, giving rise to `white dust` at the surface of the print after storage."}
{"text": "1. Technical Field\nThis invention relates to knee braces more particularly to a securing apparatus to immobilize the knee joint of a human. In therapeutic treatment of the knee which includes surgery and joint replacement it is often required to stabilize the knee there after to unload the effective ligaments to impart medial and lateral stability to the knee for a precise period of time for healing.\n2. Description of Prior Art\nPrior art devices of this type have been used to stabilize and immobilize the knee are directed to a variety of knee brace configurations which are secured to the leg above and below the knee with a mechanical support linkage there between, see for example U.S. Pat. Nos. 4,407,276, 5,135,469, 5,823,931, 5,873,847, 6,461,318.\nAn articulated knee brace is illustrated in U.S. Pat. No. 4,407,276 illustrating an upper and lower leg engagement cuff with bi-lateral interconnecting elastic strap and a knee pad extending over the expose knee cup.\nU.S. Pat. No. 5,135,469 is directed to a post surgical knee brace which incrementally is adjustable. The brace has pairs of upper and lower struts interconnected by a pair of oppositely disposed ratio-swing hinges which allow for incremental angularly bi-lateral adjustments.\nA knee brace is described in U.S. Pat. No. 5,823,931 in which an upper and lower leg engagement arms are secured to the leg above and below the knee. A bisymmetric hinge interconnects the arms and has a pair of attachment levers extending there from, each with its own pivot access in vertically aligned and spaced relation to one another.\nU.S. Pat. No. 6,561,318 claims an anatomical brace with rapid release securement members for protecting a uniform pivoting joint is disposed between the respective limb engagement members.\nAn orthopedic splint can be seen in U.S. Pat. No. 5,873,847 disclosing a pair of end plates interconnected by continuous adjustable straps to limit flexing and extension with a spring to pre-load same.\nExhibit A and B are set forth below illustrating catalog sheets which refer to the prior art disclosure.\nExhibit A is a catalog sheet from Alimed, Inc., page BQ53 illustrates a number of different immobilizer/and stabilization devices for knees having a variety of posterior medial and lateral stops with contoured leg wraps and adjustable bi-lateral hinge assemblies.\nExhibit B is a catalog sheet from PEL illustrating a number of default knee immobilizers utilizing wrap around fabric panels with interdisposed splints and deformable support stays as an example of the currently available immobilizations in the market place."}
{"text": "This section is intended to provide a background or context to the invention that is recited in the claims. The description herein may include concepts that could be pursued, but are not necessarily ones that have been previously conceived or pursued. Therefore, unless otherwise indicated herein, what is described in this section is not prior art to the description and claims in this application and is not admitted to be prior art by inclusion in this section.\nIn recent years digital imaging cameras have been commonly used as imaging devices for photographing subjects. A conventional digital camera (such as a digital still camera or digital video camera) generally acquires an image by using an imaging element or the like, and records the image as digital image data in an internal memory or integrated circuit card provided inside the camera. In an imaging device such as the above-described digital imaging camera, a lens is provided for the purpose of focusing a subject image onto the imaging element with the subject is being photographed. By controlling a focusing of the lens, a focusing distance can be aligned with the subject being photographed.\nTo address a wide spectrum of customer needs, mobile phones are rapidly converging into multi-purpose devices that incorporate multiple different products including digital imaging cameras. Digital imaging cameras nowadays are often seen as an integrated part of any mobile phone. However, the digital imaging features of these camera phones are mostly less functional and reliable than their counterparts, stand-alone digital imaging cameras. To address this gap of functionality and reliability, camera phones manufacturers strive to offer more sophisticated features and technological know-how to ensure high quality photo images with the digital imaging features of their phones.\nOne important feature to be included in a camera system is autofocus (AF). The AF feature is designed to allow a user of the camera system to obtain a correct focus of a subject of particular interest without manually adjusting the lens focal length.\nAll automatic focusing algorithms can be divided into two broad categories: active and passive. While active AF can be achieved using the external sensors (e.g. infrared), many consumer-level digital and cell phone cameras utilize passive AF. Passive AF is used to determine the correct focus of an image by using contrast or image sharpness measurements of the image. It can be understood that these measurements for AF may be performed for one or several respective areas or windows of the image.\nGenerally, a passive AF is performed utilizing an AF algorithm to calculate a measurement of image sharpness or focus corresponding to the image in order to determine a best in-focus setting. The in-focus setting is usually found by choosing a respective lens position where the image sharpness or focus is maximal. Then the determined in-focus setting can be used for an AF of the image.\nIt is noted that the operation of the AF techniques of the prior art digital cameras may deviate in their details somewhat from the description given above. However, it can be realized that prior art AF systems, such as the ones mentioned above, may be subject to certain limitations. These limitations can become apparent for certain cases, such as for cases where a resulting calculation of an image sharpness or focus does not correspond to a particular setting of the image, or where a condition of the subject or the camera interferes with the abilities of the AF system. Some of these limitations are described below in non-limiting terms.\nAccording to the numerous surveys, a highly frequent use of cameras in mobile devices is portrait imaging. Portrait imaging may be seen as the photographing of a single person or a group of people at a shorter distance, such as 1.5-3 meters, set to a more distant background which may constitute contrasting scenery. In a portrait image the single person or group of people of the image may be smaller parts of the image in comparison to the scenery. As such, a problem can exist where an AF technique, such as the passive AF technique as described above, may not choose the best in-focus position desired for the photograph.\nIn addition, it is common that while a photograph is being taken the image lighting may be low which may tend, among other things, to render parts of the image less distinguishable from other parts of the image, thus effecting an AF. Further, it is a common occurrence that during an AF the camera may inadvertently move during a photograph due to a shaking of a hand holding the camera and/or a movement of the subject being photographed. These common occurrences can subsequently lead to inaccuracies in an AF process.\nExemplary embodiments of the invention address potential AF problems, such as the ones mentioned above by example."}
{"text": "Excavating tooth assemblies provided on digging equipment such as excavating buckets or the like typically comprise a relatively massive adapter portion which is suitably anchored to the forward bucket lip and has a reduced cross-section, forwardly projecting nose portion, and a replaceable tooth point having formed through a rear end thereof a pocket opening that releasably receives the adapter nose. To captively retain the point on the adapter nose, aligned transverse openings are formed through these interengageable elements adjacent the rear end of the point, and a device commonly referred to as a flex pin or locking pin is driven into these openings.\nWhile locking pins have a variety of configurations, a widely used version, as representatively illustrated in U.S. Pat. No. 3,526,049 to Nichols and U.S. Pat. No. 3,685,178 to Ratkowski, typically comprises elongated, straight metal locking and wedge members which are laterally spaced apart and intersecured by an elongated central elastomeric element. As the locking pin is being driven into the aligned point and adapter nose openings, the elastomeric element is compressed and, when the pin is driven to its installed position, laterally urges a detent portion formed on one of the two metal portions of the point into engagement with a suitably configured portion of the adapter nose to captively retain the flex pin within the point and adapter openings. With the flex pin in its operative position within such openings, the elastomeric element is in a state of partial compression, rear surfaces of the tooth point openings bear against opposite end portions of the locking member, and a forward surface of the adapter nose opening bears against a longitudinally central portion of the wedge member. Forwardly directed tooth point removal forces encountered during the excavating process cause the tooth point to be driven forwardly relative to the adapter to thereby move the locking member closer to the elastomeric element, the opposite ends of the locking member preventing forward removal of the tooth point.\nTwo primary problems and disadvantages are present in this type of conventional flex pin construction--each of which is related to failure of the central elastomeric element. First, as the flex pin is being driven into the aligned tooth point and adapter nose openings the locking and wedge members tend to be moved longitudinally relative to one another. Thus, if the driving-in process is not carefully performed, this relative longitudinal movement can easily shear the elastomeric element, thereby mining the flex pin. Secondly, excessive forwardly directed tooth point removal loads can laterally move the locking member close enough to the wedge member to overcompress and thereby split the elastomeric element.\nVarious attempts have previously been made to better protect the critical central elastomeric portion of the flex pin by altering the essentially straight configuration of the locking and wedge member portions utilized in flex pin structures such as those depicted in the Nichols and Ratkowski patents. One such proposed solution, as exemplified in U.S. Pat. No. 4,192,089 to Schwappach and U.S. Pat. No. 4,446,638 to Novotny et al., is to form a central lateral recess in a front side portion of the locking member and to shorten the wedge member so that it is laterally movable into such recess against the resilient force of the central elastomeric element. With the elastomeric element in an uncompressed condition the opposite ends of the wedge member underlie the opposite end surfaces of the recess so that as the flex pin is being driven into the point and adapter openings one of the wedge member ends is driven into engagement with its adjacent recess end surface. This limits the relative longitudinal travel between the locking and wedge members to thereby limit the shear stress imposed upon the elastomeric element.\nIn an attempt to similarly limit the lateral compressive stress imposed on the elastomeric element, the maximum distance which the wedge member may be laterally moved into the locking member recess is limited to a distance less than the front-to-rear thickness of the elastomeric element by causing opposite end portions of the wedge member to rigidly engage portions of the locking member during travel of the wedge member into the locking member recess. In the Schwappach patent this inward travel limitation is achieved by forming on the opposite wedge member ends rearwardly directed projections which are engageable with the rear side surface of the locking recess. In the Novotny et al patent a similar result is achieved by forming forwardly facing shoulders posited adjacent opposite ends of the recess which are adapted to rigidly engage opposite end portions of the wedge member during its lateral travel into the recess. Somewhat similar schemes for protecting elastomeric flex pin portions are evidenced in U.S. Pat. No. 2,927,387 to Drover and U.S. Pat. No. 3,126,654 to Eyolfson et at.\nWhile this conventional method of limiting lateral compression of the elastomeric element represents an improvement over somewhat simpler flex pin structures such as those depicted in the Nichols and Ratkowski patents, it creates significant structural problems in the wedge member. Specifically, when the wedge member is moved to its \"stopped\" position within the locking member recess a large rigid bending load is imposed thereon by the forward surface of the adapter nose opening which bears against a central rear side portion of the wedge member. To adequately strengthen the wedge member against such bending load it is necessary to appropriately increase its front-to-rear thickness. This thickening, in turn, typically requires that undesirable design modifications be made to one or all of the elastomeric elements, the locking member and the adapter nose opening.\nSpecifically, it is well known that the overall strength of an adapter nose is, generally speaking, inversely proportional to the size of the flex pin opening formed therethrough. Thus, if it is desired to maintain a given front-to-rear length of the adapter nose opening, the necessary thickening of the wedge member requires that the front-to-rear thickness of one or both of the elastomeric element and the locking member be correspondingly reduced. Reducing the thickness of the locking member, of course, structurally weakens the flex pin, while reducing the thickness of the elastomeric element reduces the resiliency of the flex pin and the potential lateral travel between its rigid elements. Of course, neither of these results is desirable.\nIf, on the other hand, the front-to-rear thickness of the elastomeric element and the locking member are maintained, the thickening of the wedge member requires that the front-to-rear length of the adapter nose opening be correspondingly increased. This, of course, undesirably weakens the adapter nose.\nTherefore, a need exists for a locking pin which eliminates the use of an elastomeric element altogether. Such a locking pin would not experience the problems of dimensional limitations due to the thickness of the elastomeric element. Nor would it be limited to environments safe for elastomeric materials. A need exists for a one-piece locking pin, thereby eliminating the need to store various elements at the job site. A one-piece design would also limit the risk of error in installing the locking pin."}
{"text": "1. Technical Field\nThe present invention relates to an illumination device and a projector. In particular, the present invention relates to an illumination device used for a projector.\n2. Related Art\nA liquid crystal display device or a digital micromirror device (DMD) that is used for a spatial light modulation device in a projector has a characteristic in that luminance of an image is substantially regulated during one frame period of an image signal.\nIn the case in which a hold-type spatial light modulation device is used, a motion image response may be degraded due to a motion image blurring occurring at the time at when motion images are displayed.\nIn order to desirably obtain the motion image response, an impulse-type display method can be applied, in that a display light is generated during a period sufficiently shorter than the one frame period.\nIn the case of the hold-type spatial light modulation device, the impulse-type display method is performed by controlling the spatial light modulation device or by intermittent illumination that intermittently applies illumination light.\nIn such impulse-type display methods, the impulse-type display method due to the intermittent illumination is excellent in that the efficiency of illumination is not degraded and the dynamic range of the spatial light modulation device is not narrow.\nConventionally, as disclosed in, for example, Japanese Unexamined Patent Application, First Publication No. 2000-275604 and Japanese Unexamined Patent Application, First Publication No. 2004-325576, the technique in which the impulse-type display method is performed by the intermittent illumination of the spatial light modulation device is suggested.\nIn the technique in Japanese Unexamined Patent Application, First Publication No. 2000-275604, back lights formed on each illumination region are controlled so as to illuminate or turn off.\nIn this technique, there is an advantage in that it is possible to improve the motion image response due to controlling the illumination light.\nIn the case of using back lights, it is difficult to adjust each of the back lights in which the dispersion of the amount of illumination or deterioration of illumination with time occur.\nIn an optical system including a light source and a spatial light modulation device, a special expansion in a space in which a luminous flux exists is shown by the product of the emission area and radiation angle (etendue, Geometrical Extent).\nSince the etendue is restricted, it is difficult to efficiently utilize the light emitted from the light source depending on increasing the emission area due to arraying the light sources.\nIn Japanese Unexamined Patent Application, First Publication No. 2004-325576, the illumination light is scanned onto a projection surface of the spatial light modulation device.\nIn the impulse-type display method, though the motion image response can be improved, a flicker of the images on the display is easily generated.\nThe flicker is easily noticed by a viewer, especially, in the case of displaying a still image on the display.\nIn Japanese Unexamined Patent Application, First Publication No. 2004-325576, since it is impossible to select scanning or non-scanning in accordance with images, it is difficult to suppress the occurrence of flicker.\nAs described above, though a high efficiency of illumination and a desirable motion image response can be obtained, there is the problem in that it is difficult to suppress the occurrence of flicker."}
{"text": "1. Field\nThe present application relates generally to wireless communications, and more specifically to systems, methods, and devices having wireless dual concurrent bands (DCB) or multiple concurrent bands (MCB), including operations achieving higher overall throughput and improving communication efficiency via band and/or channel switching and buffer management.\n2. Background\nIn many telecommunication systems, communications networks are used to exchange messages among several interacting spatially-separated devices. Networks may be classified according to geographic scope, which could be, for example, a metropolitan area, a local area, or a personal area. Such networks may be designated respectively as a wide area network (WAN), metropolitan area network (MAN), local area network (LAN), or personal area network (PAN). Networks also differ according to the switching/routing technique used to interconnect the various network nodes and devices (e.g., circuit switching vs. packet switching), the type of physical media employed for transmission (e.g., wired vs. wireless), and the set of communication protocols used, e.g., Internet protocol suite, Synchronous Optical Networking (SONET), Ethernet, etc.\nWireless networks are often preferred when the network elements are mobiles and thus have dynamic connectivity needs, or if the network architecture is formed in an ad hoc, rather than fixed, topology. Wireless networks may employ intangible physical media in an unguided propagation mode using electromagnetic waves in the radio, microwave, infra-red, optical, etc. frequency bands. Wireless networks advantageously facilitate user mobility and rapid field deployment when compared to fixed wired networks.\nThe devices in a wireless network may transmit/receive information between each other. The information may include packets, which in some aspects may be referred to as data units. The packets may include overhead information (e.g., header information, packet properties, etc.) that helps in routing the packet through the network, identifying the data in the packet, processing the packet, etc., as well as data, for example user data, multimedia content, etc. as might be carried in a payload of the packet.\nRecent evolutions of IEEE 802.11 standard expand Wi-Fi onto multiple frequency bands. These bands include: 2 GHz frequency band for IEEE 802.11b/g/n, 5 GHz frequency band for IEEE 802.11a/n/ac, 60 GHz frequency band for IEEE 802.11ad, 900 MHz frequency band for IEEE 802.11ah and TVWS band for IEEE 802.11af. Recent technology enables a wireless device to concurrently operate on multiple bands. With a MCB station (STA) and a MCB access point (AP), multiple links may be established between the STA and the AP. Many Wi-Fi products support DCB, usually including the 2 GHz frequency band and 5 GHz frequency band. A DCB STA and a DCB AP may operate on both frequency bands concurrently and achieve higher overall throughput. However, such a DCB operation mode may cause higher power consumption, and therefore may not be beneficial. Therefore, it would be desirable if a DCB STA or a MCB STA can switch an operation frequency band with an associated AP dynamically."}
{"text": "The steroid hormone cortisol is a key regulator of many physiological processes. However, an excess of cortisol, as occurs in Cushing's Disease, provokes severe metabolic abnormalities including: type 2 diabetes, cardiovascular disease, obesity, and osteoporosis. Many patients with these diseases, however, do not show significant increases in plasma cortisol levels. In addition to plasma cortisol, individual tissues can regulate their glucocorticoid tone via the in situ conversion of inactive cortisone to the active hormone cortisol. Indeed, the normally high plasma concentration of cortisone provides a ready supply of precursor for conversion to cortisol via the intracellular enzyme 11-beta-hydroxysteroid dehydrogenase type I (11beta-HSD1).\n11beta-HSD1 is a member of the short chain dehydrogenase superfamily of enzymes. By catalyzing the conversion of cortisone to cortisol, 11beta-HSD1 controls the intracellular glucocorticoid tone according to its expression and activity levels. In this manner, 11beta-HSD1 can determine the overall metabolic status of the organ. 11beta-HSD1 is expressed at high levels in the liver and at lower levels in many metabolically active tissues including the adipose, the CNS, the pancreas, and the pituitary. Taking the example of the liver, it is predicted that high levels of 11beta-HSD1 activity will stimulate gluconeogenesis and overall glucose output. Conversely, reduction of 11beta-HSD1 activity will downregulate gluconeogenesis resulting in lower plasma glucose levels.\nVarious studies have been conducted that support this hypothesis. For example, transgenic mice expressing 2× the normal level of 11beta-HSD1 in only the adipose tissue show abdominal obesity, hyperglycemia, and insulin resistance. (H. Masuzaki, J. Paterson, H. Shinyama, N. M. Morton, J. J. Mullins, J. R. Seckl, J. S. Flier, “A Transgenic Model of Visceral Obesity and the Metabolic Syndrome”, Science, 294:2166-2170 (2001). Conversely, when the 11beta-HSD1 gene is ablated by homologous recombination, the resulting mice are resistant to diet induced obesity and the accompanying dysregulation of glucose metabolism (N. M. Morton, J. M. Paterson, H. Masuzaki, M. C. Holmes, B. Staels, C. Fievet, B. R. Walker, J. S. Flier, J. J. Mullings, J. R. Seckl, “Novel Adipose Tissue-Mediated Resistance to Diet-induced Visceral Obesity in 11β-Hydroxysteroid Dehydrogenase Type 1-Deficient Mice”, Diabetes, 53: 931-938 (2004). In addition, treatment of genetic mouse models of obesity and diabetes (ob/ob, db/db and KKAy mice) with a specific inhibitor of 11beta-HSD1 causes a decrease in glucose output from the liver and an overall increase in insulin sensitivity (P. Alberts, C. Nilsson, G. Selen, L. O. M. Engblom, N. H. M. Edling, S. Norling, G. Klingstrom, C. Larsson, M. Forsgren, M. Ashkzari, C. E. Nilsson, M. Fiedler, E. Bergqvist, B. Ohman, E. Bjorkstrand, L. B. Abrahmsen, “Selective Inhibition of 11β-Hydroxysteroid Dehydrogenase Type I Improves Hepatic Insuling Sensuitivity in Hyperglycemic Mice Strains”, Endocrinology, 144:4755-4762 (2003)). Furthermore, inhibitors of 11beta-HSD1 have been shown to be effective in treating metabolic syndrome and atherosclerosis in high fat fed mice (Hermanowoki-Vosetka et al., J. Eg. Med., 202(4):517-527 (2002)). Based in part on these studies, it is believed that local control of cortisol levels is important in metabolic diseases in these model systems. In addition, the results of these studies also suggest that inhibition of 11beta-HSD1 will be a viable strategy for treating metabolic diseases such as type 2 diabetes, obesity, and the metabolic syndrome.\nLending further support to this idea are the results of a series of preliminary clinical studies. For example, several reports have shown that adipose tissue from obese individuals has elevated levels of 11beta-HSD1 activity. In addition, studies with carbenoxolone, a natural product derived from licorice that inhibits both 11beta-HSD1 and 11beta-HSD2 (converts cortisol to cortisone in kidney) have shown promising results. A seven day, double blind, placebo controlled, cross over study with carbenoxolone in mildly overweight individuals with type 2 diabetes showed that patients treated with the inhibitor, but not the placebo group, displayed a decrease in hepatic glucose production (R. C. Andrews, O. Rooyackers, B. R. Walker, J. Clin. Endocrinol. Metab., 88:285-291 (2003)). This observation is consistent with the inhibition of 11beta-HSD1 in the liver. The results of these preclinical and early clinical studies strongly support the concept that treatment with a potent and selective inhibitor of 11beta-HSD1 will be an efficacious therapy in patients afflicted with type 2 diabetes, obesity, and the metabolic syndrome."}
{"text": "The present invention relates to a process for producing fiber-reinforced silicon carbide composites. Specifically, the present invention relates to a process for producing fiber-reinforced silicon carbide composites which are suitable for a variety of applications requiring improved toughness, such as aerospace high-temperature structural members, gas turbine members, fusion reactor materials, furnace members, heater materials, and artificial bones.\nSilicon carbide ceramics are light-weight and are excellent in, for example, heat resistance, abrasion resistance, and corrosion resistance; they have recently come into wide use as, for instance, high-temperature corrosion-resistant members, heater members, abrasion resistant members, as well as abrasives and grindstones. Such silicon carbide ceramics are, however, low in fracture toughness, and have not yet been used in practice as structural members for use at high temperatures.\nCeramic composites compounded with fibrous reinforcements to improve the toughness of such ceramics have been intensively investigated in recent years. Such fiber-reinforced silicon carbide composites are generally manufactured by, for example, (1) an organometallic polymer impregnation pyrolysis (PIP) process, (2) a chemical vapor infiltration (CVI) process, and (3) infiltration of molten silicon (reaction sintering process).\nHowever, the organometallic polymer impregnation pyrolysis (PIP) process has the following disadvantages in practice. According to this process, only low density and low strength can be obtained by a single impregnation, and the impregnation and pyrolysis procedure must be repeated about ten times to reduce the open porosity to 10% or less to and thereby improve strength characteristics. The process therefore requires a long time for production. The chemical vapor infiltration process (2) can provide products having complicated shapes at comparatively low temperatures of about 1100xc2x0 C., but it requires a very long time, of as much as several months, for infiltration, and gases used therefore are toxic. In addition, composites having an open porosity not exceeding 5% cannot be significantly obtained by the single use of the process (1) or (2).\nIn contrast, the reaction sintering process (3) requires only a short reaction time and can yield dense composites in a short time period. According to a process employed in Deutsche Forschungs-und Versuchsanstalt fur Luft-und Raumfahrt (DLR), a carbon fiber-reinforced carbon-silicon carbide composite is produced by infiltrating molten silicon into cracks of a carbon-fiber-reinforced carbon composite (C/C composite) to convert part of the matrix carbon into silicon carbide. This process utilizes a phenomenon in which glassy carbon does not react significantly with molten silicon to avoid the reaction between the carbon fiber and the silicon. In the process, however, the mechanical properties depend to a large degree on shapes of the cracks, which shapes in turn depend, for example, on the type of the carbon fiber and on the heat treatment temperature. High mechanical properties can therefore only be provided under specific production conditions suitable for the fiber to be used. The reaction sintering process (3) also includes a process of infiltrating molten silicon into a composite of a carbon powder and a fiber, which is performed by General Electric Co., USA, and Toshiba Corporation, Japan. This process is also disadvantageous in that a large quantity of silicon must be infiltrated, which invites free silicon to remain in large amounts of about 15%, or causes a reaction between the fiber and silicon unless the fiber is coated with boron nitride (BN) or the like.\nAfter investigations on similar ceramic composites, the present inventors previously found that a unidirectional carbon-fiber-reinforced silicon carbide composites having a flexural strength of about 200 to 300 MPa even though having a large open porosity of 30% can be produced by compounding a matrix of a mixture of silicon powder and phenol resin with carbon fiber, and firing the composite in an inert atmosphere. The resultant composite contains a matrix of silicon carbide formed by reaction sintering (Japanese Patent No. 2045825). According to the reaction, however, only porous composites having a large open porosity can be obtained, as the volume of the matrix decreases about 38%. They also found that a unidirectionally carbon-fiber-reinforced silicon carbide composite having a flexural strength of about 500 to 600 MPa, even though it has a large open porosity of about 20%, can be produced by reducing the particle of the silicon powder about 5 xcexcm or below and by adding an organometallic polymer (Japanese Patent No. 2735151), and that a two-dimensional fiber-reinforced composite having a relatively large open porosity of about 15% as a two-dimensional fiber-reinforced silicon-carbide carbon composite, but having a flexural strength of about 300 MPa, can be obtained by heat-treating a green body containing a fibrous woven fabric reinforcement at a temperature at which silicon does not react with carbon, repeating impregnation and carbonization of a phenol resin, and finally forming silicon carbide (Japanese Patent No. 2879675).\nIt is an object of the present invention to solve the various problems in the production of fiber-reinforced silicon carbide composites according to the conventional silicon melt infiltration technique, and to provide a process for producing fiber-reinforced silicon carbide composites, which process can easily yield a composite having a high toughness even if they are to have complicated shapes.\nSpecifically the invention is directed to provide a process for easily producing fiber-reinforced silicon carbide composites even if they are to have complicated shapes, by covering a fiber bundle region with glassy carbon derived from resin, forming a porous region in a specific portion of a matrix by a reaction between silicon powder and carbon derived from the resin to form silicon carbide, which reaction is accompanied by volumetric reduction, and subjecting the porous region to infiltration of molten silicon. The resultant fiber-reinforced silicon carbide composite is high in toughness and its strength is not deteriorated, without coating the fiber surface with, for example, BN.\nAfter intensive investigations on the production of fiber-reinforced silicon carbide composites to achieve the above objects, the inventors found that a dense fiber-reinforced silicon carbide composite showing nonlinear fracture can be obtained by preparing and molding a prepreg including silicon powder, carbon source resin and fiber to yield a green body, or laminating a prepreg containing resin with a prepreg containing silicon powder and resin in alternate order and molding the laminate to yield a green body; carbonizing the molded green body at a temperature of about 900xc2x0 C. to 1350xc2x0 C. in an inert atmosphere; preferably, impregnating the carbonized composite with a resin and carbonizing the impregnated composite at a temperature of about 900xc2x0 C. to 1350xc2x0 C. in an inert atmosphere, and repeating this impregnation-carbonization procedure; subjecting the composite to reaction sintering at a temperature of about 1300xc2x0 C. or higher in vacuo or in an inert atmosphere, and finally infiltrating molten silicon into the sintered composite at a temperature of about 1300xc2x0 C. to 1800xc2x0 C. in vacuo or in an inert atmosphere. The present invention has been accomplished on the basis of the above findings.\nSpecifically, in the process for producing a fiber-reinforced silicon carbide composite of the invention, (i) a mixture of silicon powder, a carbon source resin and reinforcement fiber is carbonized in an inert atmosphere, and the resultant carbonized composite is impregnated with resin and is subjected to carbonization or the like, or (ii) a mixture of (a) a matrix containing silicon powder and carbon source resin, and porous fiber as a container, and (b) reinforcing fiber containing carbon source resin is carbonized in an inert atmosphere, and the carbonized composite is impregnated with a resin and subjected to carbonization or the like. By the procedure of (i) or (ii), carbon in a fiber bundle region as a reinforcement is densified to prevent the fiber from reacting with molten silicon, then open pores are formed in a matrix region by reaction sintering, and molten silicon is infiltrated into only the open pores in the matrix in vacuo or in an inert atmosphere.\nAccording to the inventive process, dense fiber-reinforced silicon carbide composites can be easily obtained without deteriorating the strength of the fiber, even if they are of complicated shape.\nPreferred carbon source resins to be used in the invented process include, but are not limited to, phenol resins, furan resins, pitch, and polycarbosilane and other organometallic polymers. Each of these resins may be used alone or in combination. The resin may further include, for instance, a carbon powder, a graphite powder, carbon black, or an aggregate; and/or an antioxidant such as silicon carbide, silicon nitride, mullite, molybdenum disilicide, molybdenum, boron carbide, or boron powder.\nThe silicon may be either pure metallic silicon or a silicon alloy with, for example, magnesium, aluminum, titanium, chromium, manganese, iron, cobalt, nickel, copper, zinc, zirconium, niobium, or molybdenum.\nIn addition, a fibrous woven fabric reinforcement is employed in the invented process, and preferred fibrous woven fabrics to constitute the fibrous reinforcement include, but are not limited to, woven fabrics of carbon fibers, silicon carbide fibers, silicon nitride fibers, and oxide fibers. The carbon fibers may be either of pitch fibers or acrylonitrile fibers. Each of these fibrous woven fabrics can be used independently or in combination. Nonwoven fabrics or laminates of transversely laid sheet-like unidirectional fiber prepregs may also be used instead of fibrous woven fabrics.\nAs the fiber in the matrix to contain the silicon powder and resin, porous woven or nonwoven fabrics are preferred composed of, for example, carbon fibers, silicon carbide fibers, silicon nitride fibers or oxide fibers, as well as porous sheets composed of carbon, silicon carbide, a resin or a plastic. The use of the container can be omitted by coating a resin-impregnated woven fabric prepreg with a mixture of silicon powder and carbon source resin.\nAccording to the inventive process as described above, dense fiber-reinforced silicon carbide composites can easily be produced by covering a fiber bundle region of a two-dimensional fiber reinforced silicon carbide composite with amorphous carbon derived from an impregnated resin to densify the region; forming open pores in spaces between the fibers or between the fibrous woven fabrics with the use of reaction sintering of, for example, a phenol resin with silicon powder; and infiltrating molten silicon into the open pores.\nIn addition, and advantageously, the invented process can easily yield green bodies (molded parts) having complicated shapes because hot pressing or other techniques are not required. Furthermore, the process can yield dense composites without deteriorating the fiber, as molten silicon is infiltrated in the last step.\nNext, preferred embodiments of the present invention will be described.\nAccording to the first embodiment of the invented process (hereinafter simply referred to as xe2x80x9cthe first processxe2x80x9d), a continuous woven fabric as a fibrous reinforcement is impregnated with dissolved resin, is dried, and is further impregnated with a slurry of a mixture of resin and a silicon powder; is dried to give a prepreg; and the prepreg is molded into a desired shape. The molding is carried out by, for example, dissolving resin in suitable solvent, impregnating a continuous fiber woven fabric with the resin solution, removing the solvent by drying, separately preparing a slurry composed of resin, silicon powder and solvent, impregnating continuous fiber woven fabric with the slurry, removing the solvent by drying, molding the dried product in an appropriate die at a temperature of, generally, about 100xc2x0 C. to 250xc2x0 C. into a desired shape. In this procedure, the dissolved liquid resin is incorporated into spaces between fiber bundles, but the powdered silicon, which generally has a relatively large particle size of about 10 xcexcm, is unevenly distributed and present only in spaces between the fibrous woven fabric.\nIn addition, and advantageously, the following process according to the second embodiment (hereinafter simply referred to as xe2x80x9cthe second processxe2x80x9d) is also effective. According to the second process, continuous fiber woven fabric as a fibrous reinforcement is impregnated with dissolved resin, and is dried to give fiber prepregs; separately, a slurry is prepared by mixing resin and silicon powder, nonwoven fabric is impregnated with the slurry, and is dried to give matrix prepregs; and the fiber prepregs and the matrix prepregs are laminated in alternate order, and the laminate is molded into a desired shape. The laminate may have each layer comprising a plurality of individual prepregs. The molding can be performed by, for example, dissolving a resin in a suitable solvent, impregnating continuous fiber woven fabric with the resin solution, removing the solvent by drying to give fiber prepregs, separately preparing a slurry from resin, silicon powder and a solvent, impregnating nonwoven fabric with the slurry, removing the solvent by drying to give matrix prepregs, laminating an each an adequate number of these prepregs in alternate order, and molding the laminate in an appropriate die at a temperature of, generally, about 100xc2x0 C. to 250xc2x0 C. into a desired shape.\nThe silicon powder to be used in the invented process is preferably fine powder, and particularly fine powder having a mean particle size not exceeding 20 xcexcm. Silicon powders having a relatively large diameter should preferably be pulverized with, for example, a ball mill, before use.\nThe obtained green body (molded body) is then carbonized at a temperature of about 900xc2x0 C. to 1350xc2x0 C. in an argon or other inert atmosphere to avoid the reaction of carbon derived from the resin with silicon. The carbonized composite has regions composed of the resin alone, and regions composed of a mixture of the silicon powder and carbon derived from the resin, but it has a high porosity, because the resin is pyrolyzed into carbon.\nThe fired body is impregnated with a resin solution in vacuo or under pressure, and is carbonized at a temperature of about 900xc2x0 C. to 1350xc2x0 C. in an argon or other inert atmosphere to avoid the reaction of carbon derived from the resin with silicon. A combination procedure of the resin impregnation and carbonization is repeated a predetermined number of times, preferably one to five times, and the resultant composite is fired at a temperature of 1300xc2x0 C. or higher in vacuo or in an argon or other inert atmosphere to react carbon derived from the resins with silicon to form porous silicon carbide in spaces between the fiber woven fabric.\nA composite is thus obtained comprising fiber bundles coated with dense carbon so as not to react with molten silicon, and a porous silicon carbide matrix between layers of the fiber woven fabric.\nThe composite is heated to a temperature of about 1300xc2x0 C. to 1800xc2x0 C. in vacuo or in an inert atmosphere to infiltrate molten silicon into open pores of the composite (fired body), which open pores include the porous silicon carbide matrix between the layers of the fiber woven fabric, to give a fiber reinforced silicon carbide composite. The reaction sintering of silicon with carbon and infiltration of molten silicon can be performed in the same heat treatment.\nThe proportion of the silicon powder to the resin to be used in the inventive process should preferably be selected from such a range that the atomic ratio of silicon to carbon derived from the resin, Si/C, ranges from 0.1 to 4. The prepregs may be prepared by, as described above, coating the fabric with phenol resin or another resin alone, and coating the coated fabric with a slurry containing silicon powder and phenol resin or another resin. In other words, the proportion of silicon to carbon in the composite can be distributed unevenly. The proportion of silicon carbide matrix to the fiber woven fabric reinforcement in the fiber-reinforced composite can be freely selected according to the application of the composite, but generally, the fiber occupies about 15% to 65% by volume in the composite.\nIn the invention, the prepreg containing the reinforcing fiber woven fabric can be prepared by any of the following techniques: (i) a technique of using resin from the beginning to yield a prepreg, as described above; (ii) a technique which comprises laminating a container for silicon powder and resin, or silicon powder and aggregate powder with a fiber woven fabric containing no resin to yield a green body, and impregnating the green body with resin to yield a preform; and (iii) a technique which comprises laminating a container for neither silicon nor resin with reinforcing fiber woven fabric containing resin to yield a green body, and impregnating the green body with a slurry containing silicon powder and resin.\nAnother resin impregnation can be performed after the reaction sintering of carbon from phenol resin or other resin with silicon powder, which is completed at temperatures of 1300xc2x0 C. or higher. Separately, the reaction sintering can be conducted prior to, or concurrently with, the infiltration of molten silicon."}
{"text": "In recent years, we have greeted the age of multimedia in which audio, video, and other data are integrally handled, and the conventional information media, i.e., means for transmitting information between persons, such as newspapers, magazines, televisions, radios, and telephones, have been grasped as subjects of multimedia. Generally, “multimedia” does not mean only representing characters, but means representing especially images, simultaneously in relation with diagrams, speech, and the like. In order to adopt conventional information as information media as subjects of multimedia, it is required to represent the information in a digital form.\nWhen the information quantity possessed by the respective information media described above are estimated as digital information quantity, the information quantity per character is 1-2 bytes, while the information quantity of more than 64 kbits per second (telecommunication quality) is required for audio and the information quantity of more than 100 Mbits per second (current television broadcasting quality) is required for moving picture. Therefore, it is not practical to handle such tremendous amounts of data of the above-described information media as they are in digital formats. For example, though visual phones have already been put to practical use by ISDN (Integrated Services Digital Network) having a transmission rate of 64 kbps˜1.5 Mbps, it is impossible to transmit images of television cameras as they are by the ISDN.\nIn these circumstances, information compression techniques are required. For example, in the case of visual telephones, the moving picture compression techniques standardized as H.261 and H.263 by ITU-T (International Telecommunication Union-Telecommunication Sector) are employed. Further, according to the information compression technique of MPEG1, it is possible to contain video information together with audio information in an ordinary music CD (compact disk).\nMPEG (Moving Picture Experts Group) is an international standard of data compression for a moving picture, i.e., pixel values of a moving picture. MPEG1 is a standard for compressing pixel values of a moving picture to 1.5 Mbps, i.e., data of a television signal to about 1/100. Further, while the transmission rate in MPEG1 is mainly limited to about 1.5 Mbps, in MPEG2 which is aimed at standardization to meet a request for a higher image quality, the limitation is relaxed with pixel values of a moving picture being compressed to 2˜15 Mbps.\nFurther, under the existing circumstances, standardization of MPEG4 has been almost completed by the working group for standardization of MPEG1 and MPEG2 (ISO/IEC JTC1/SC29/WG11), which enables coding and handling in object units and realizes new functions demanded in the multimedia age. While MPEG4 has initially aimed at standardization of a coding method of a low bit rate, the aim of the standardization thereof is now extended to a more versatile coding process of a high bit rate or an interlaced image. One of characteristics of MPEG4 is coding simultaneously plural image sequences and transmitting the same. This enables one image scene to be composed of plural images. The foreground and the background can be different image sequences, and the frame frequency, the image quality and the bit rate thereof can be individually changed. Thereby, plural images can be arranged in the horizontal or vertical direction like in a multi-screen, and it is enabled for the user to extract or enlarging-display only a desired image. It is general that only pixel values are coded for the background similarly in MPEG2, while, as for the foreground, a pixel value signal indicating pixel values of the object as well as a shape signal indicating the shape of the object are coded. Commonly, the coding of the foreground are known as coding in object units. The displayed image is a video composed of the respective decoded images.\nFIGS. 8(a)-8(f) are diagrams for explaining video composition in object units. FIG. 8(a) shows pixel values of a foreground video (hereinafter, simply referred to as foreground) fv1 of a balloon which composes a display video. FIG. 8(b) shows a shape value of a shape signal fml corresponding to the foreground fv1 in FIG. 8(a). FIG. 8(c) shows pixel values of a foreground fv2 of a human which composes the display video. FIG. 8(d) shows a shape value of a shape signal fm2 corresponding to the foreground fv2 in FIG. 8(c). FIG. 8(e) shows pixel values of a background video (hereinafter, simply referred to as background) rv. FIG. 8(f) shows the display video which is obtained by composing the foregrounds fv1 and fv2 of FIGS. 8(a) and 8(c) with the background rv of FIG. 8(e).\nIn the screens of the shape signals of FIGS. 8(b) and 8(d), parts painted black in the screens show areas in which the corresponding pixel values exist, i.e., inside an object, and white parts show areas in which no corresponding pixel values exist, i.e., outside the object. A wording that “a pixel value is significant/insignificant” is sometimes used in a sense that a pixel value exists/no pixel value exists.\nFIG. 9 is a block diagram illustrating a structure of a prior art video decoding system. In this figure, reference DeMux denotes a stream demultiplexer for demultiplexing a multiplexed stream StrM. References Dec1, Dec2 and Dec3 denote decoders for decoding video streams Str1, Str2 and Str3 which have been demultiplexed by the stream demultiplexer DeMux, respectively. Reference Comp1 denotes a video composer for composing decoded videos Dout1, Dout2 and Dout3 which have been decoded by the decoders Dec1, Dec2 and Dec3, respectively. Reference Disp denotes a display unit for displaying a video Vcom which is composed by the video composer Comp1. Reference IF denotes an object selector for the user to select an object. CPU denotes a controller for controlling the decoders Dec1, Dec2 and Dec3 in accordance with the instruction of the object selector IF.\nNext, the operation of the prior art video decoding system is described. A multiplexed stream StrM is demultiplexed by the stream demultiplexer DeMux into video streams Str1, Str2 and Str3 corresponding to videos of three objects, as well as overlap information Odr indicating the order of overlap of these videos which is sent to the controller CPU. The decoder Dec1 decodes the video stream Str1 and outputs the decoded video Dout1. Similarly, the decoders Dec2 and Dec3 decode the video streams Str2 and Str3, and output the decoded videos Dout2 and Dout3, respectively. The video composer Comp1 composes these decoded videos Dout1, Dout2 and Dout3 to provide a composed video Vcom, and displays the composed video on the display unit Disp.\nOn the other hand, in the case of object unit coding, the user can switch the display/non-display of videos in object units. The user selects the display or non-display of each object by means of the object selector IF. The object selector IF notifies the video composer Comp1 of object non-display information Dsel according to this selection, and the video composer Comp1 composes only videos of objects which should be displayed, to display the composed video.\nWith using the shape value of an object, it can be judged whether a position is inside or outside the object. Accordingly, by executing an operation of selecting a button Bn within the screen using a pointer Pr as shown in FIG. 10, it is possible that the user designates a specific position in the screen and obtains information as to whether the position is inside or outside an object. To be specific, the user moves the pointer by the object selector IF and selects a button, thereby notifying the controller CPU of object selection information Psel which indicates a designated pixel position. The controller CPU makes an inquiry to the decoders Dec1, Dec2 and Dec3 corresponding to the respective objects with object judge commands Q1, Q2 and Q3 about whether the designated position is inside or outside an object such as a button on the screen. The decoders Dec1, Dec2 and Dec3 report the controller CPU whether the inquired position is inside or outside the object by object judge results A1, A2 and A3, respectively, and then the controller CPU notifies the user or applications of the object judge results A1, A2 and A3 collectively as object judge result Req.\nThe block diagram of FIG. 9 illustrating the video decoding system shows an example where three videos are composed to obtain a composed video Vcom, while the number of videos to be composed can be less than three, or more than three. In addition, in this example, one decoder is provided for each video stream, while when plural video streams can be decoded by one video decoder by time division or the like, the number of video decoders can be properly reduced.\nFIG. 11 is a block diagram illustrating a structure of the decoder Dec of the video decoding apparatus in the prior art video decoding system shown in FIG. 9. In this figure, a video stream Str, an object judge command Q, an object judge result A and a decoded video Dout correspond to one of the video streams Str1, Str2 and Str3, the object judge commands Q1, Q2 and Q3, the object judge results A1, A2 and A3, and the decoded videos Dout1, Dout2 and Dout3 in FIG. 9, respectively. Reference DecU denotes a video decoding unit for decoding the video stream Str. References MEM1, MEM2, MEM3 and MEM4 denote memories for containing decoded videos mem1, mem2, mem3 and mem4, respectively.\nNext, the operation of the decoder Dec is described. In FIG. 11, the video decoding unit DecU decodes the video stream Str and stores the decoded videos mem1, mem2, mem3 and mem4 which are obtained by the decoding, into the memories MEM1, MEM2, MEM3 and MEM4, respectively. At this time, when the video stream Str has been inter-frame coded, the decoded videos mem1, mem2, mem3 and mem4 are read from the memories MEM1, MEM2, MEM3 and MEM4 to utilize the same as reference videos at the motion compensation. Since a signal of a video having the shape is composed of three components indicating the color (YUV, RGB or the like) and a shape value A, i.e., four components in total, the four individual memories are provided to correspond to the four components, respectively. However, when practically packaged, these can be integrated in one memory.\nIn the example shown in FIG. 11, the memories MEM1, MEM2, MEM3 and MEM4 contain a luminance pixel value Yimg, two color difference pixel values Uimg and Vimg, and a shape value Aimg, respectively. Since the shape value Aimg is stored in the memory MEM4, when a pixel position such as a position pointed by the pointer is indicated by the object judge command Q from the controller CPU in FIG. 9, the memory MEM4 judges whether that position is inside or outside the object, and outputs the result as the object judge result A. The decoded videos stored in the memories MEM1, MEM2, MEM3, and MEM4 are read as pixel decoded videos Yimg, Uimg and Vimg, and a shape decoded video Aimg at a timing of display, to obtain pixel decoded videos Yout, Uout and Vout, and a shape decoded video Aout, respectively. The decoded video Dout is obtained by combining the pixel decoded videos Yout, Uout and Vout and the shape decoded video Aout.\nFIG. 12 is a diagram showing an internal structure of the memory MEM4 which contains the shape value. In this figure, reference MEM41 denotes a shape signal storage memory for containing the shape signal. Reference CMP denotes a pointed position comparison means for comparing and judging whether the pointer operated by the object selector IF in FIG. 9 points inside or outside of an object, such as a button.\nNext, the operation of the memory MEM4 is described. In FIG. 12, the shape signal storage memory MEM41 contains the shape signal mem4 which has been decoded by the decoding unit DecU in FIG. 11 as a bitmap. The pointed position comparison means CMP converts pointed position information of the pointer operated by the object selector IF in FIG. 9, which is transmitted in accordance with the object judge command Q issued by the controller CPU in FIG. 11, together with this object judge command Q, into an address of the shape signal storage memory MEM41, and judges whether or not the bitmap of the shape signal exists at that address, thereby judging whether the position pointed by the pointer is inside or outside of the video of the object, such as a button. Then, the pointed position comparison means CMP outputs the judge result to the controller CPU in FIG. 9 as the object judge result A.\nAs described above, the information as to whether a position is inside or outside an object can be obtained by using the shape value of the object. By utilizing these workings, the shape value can be used as a GUI (Graphic User Interface) operable button whose shape varies. This is what is called a “hot spot”, and, for example, when a certain position on the screen of a terminal (for example, assuming that this is a video in a form of a button) is clicked to make the terminal execute a special processing, a shape signal indicating that position is transmitted as a moving picture, thereby changing the position or shape of the button. For that purpose, not only videos including both of pixel values and shape values, but also videos having only shape values are effective, and accordingly the coding of only shape values can be also used in MPEG4.\nFor example, when FIG. 8(f) is used as background, only shape values of FIGS. 8(b) and 8(d) are coded, and a position is pointed in FIG. 8(f) by a pointer, the judgement as to whether that position is inside or outside the human or balloon can be made. Therefore, when only the judgement as to whether a position is inside or outside the video is to be made, the coding of individual pixel values of FIGS. 8(a), 8(c) and 8(e) is not required, and only the coding of the pixel value of FIG. 8(f) is required. Therefore, the coding/decoding process can be simplified, and in some cases the compression rate can be also increased by reduction of the number of pixels to be coded.\nFrom the above descriptions, it can be seen that there are three kinds of the stream of color videos (moving pictures) in MPEG4, i.e., only YUV (color signals) in the case of videos whose shapes do not vary, YUV+A (a shape signal is added to the color signals) in the case where coding in object units is carried out, and only A (shape signal) in the case of judgement as to whether a position is inside or outside a video is made.\nThere are some cases where videos are transmitted according to MPEG1 and a shape signal according to MEPG4 is added thereto. Further, there is also a case where a texture is pasted on a shape signal transmitted according to MPEG4, whereby MPEG4 and CG (Computer Graphics) are combined to display color videos.\nFIG. 13 is a diagram schematically showing the format of a video stream Str corresponding to a video of an object. In FIG. 13, reference HD denotes the header of the entire stream. References DA1, . . . , and DAX denote data of one screen, respectively. References HD1, . . . , and HDX denote headers corresponding to the screen data DA1, . . . , and DAX, respectively.\nIn the header HD of the entire stream, the video size (when the size of the video does not vary with frames) and the coding method (quantization method or information relating to arrangement of data) as well as information indicating a target which is being coded (above-mentioned YUV, YUV+A, A or the like), are coded and stored.\nIn the headers HD1, . . . , and HDX corresponding to the screen data DA1, . . . , and DAX, respectively, parameters required for the decoding, information indicating which frame is the corresponding video data or which is of I frame and P frame the video data, and the like are coded and stored.\nFIG. 14 is a diagram showing a structure of a multiplexed stream StrM which is obtained by multiplexing plural video streams Str each corresponding to a video of an object. In the example shown in FIG. 14, the video streams Str are time-divided multiplexed frame by frame, and a header MHD including the overlap information Odr is arranged between the video streams Str.\nAs described above, it is useful to code only shape values, while when a stream including no pixel value but having only the shape value is received, what becomes the pixel value which is obtained by decoding this stream is not decided in the MPEG4 standard at the present time.\nOriginally, a stream having only the shape value is created provided that it is not displayed on the receiving end. However, since the measures to be taken when this is received and decoded are not defined in the MPEG4 standard, the stream having only the shape value should not be displayed, in accordance with proposals on the application side which utilizes video communication to provide information terminals with various kinds of services.\nHowever, in many cases, video decoders created for general purposes are generally used for various applications to reduce developing costs, and these video decoders are designed to always decode and display transmitted information. Therefore, also when receiving a stream having only the shape value, the video decoder decodes this stream, and some pixel values which cannot be predicted are displayed due to that decoding, thereby giving wrong or unpleasant feelings to persons who watch the screen."}
{"text": "Semiconductor integrated circuit devices are designed using various computer aided means to produce a logical circuit in which various functional elements are identified and interconnected to form particular logic functions. Computer-aided design (CAD) or other design tools are typically used to generate lithography data based on the design and the lithography data is used to form a photomask set that is used in the actual fabrication of the semiconductor integrated circuit (IC) devices.\nA standard cell type semiconductor IC device is an application specific integrated circuit (ASIC) device designed using multiple standard cells. Each standard cell may include a group of transistor and interconnect structures that provide a boolean logic function (e.g., AND, OR, XOR, XNOR, inverters) or a storage function (flipflop or latch, for example). In a standard cell system, the standard cells are advantageously maintained in a library and various integrated circuit devices can be designed by selecting desired standard cells from the library and combining them in various arrangements to form a most suitable cell arrangement pattern that forms part of the layout of the integrated circuit device.\nThe design process for standard cell semiconductor IC devices includes constructing the integrated circuit design out of the selected standard cells that are connected together electrically, i.e. routed, using wire interconnects which are also used to route the cells to input and output terminals. The standard cells and connections between them may be stored in databases called “netlists” (i.e., lists of symbolic interconnections). The netlist describes the connectivity of the standard cells within the semiconductor IC device design. One or more computer-aided design (CAD) tools may be used to generate the netlist of the selected standard cells and the interconnections between the cells and input/output. The netlist may be used by a floor planner, placement tool, or other design tool, to place the selected cells at particular locations in the layout of a standard cell integrated circuit device.\nThe design of the standard cell semiconductor IC device may be carried out by an APR (automatic placement and routing) design tool that includes a placer and a router, by selecting standard cells from the library of standard cells and placing and routing the cells according to design instructions provided to the APR tool such as in the netlist. The placer determines the optimum location of each standard cell of the integrated circuit on the semiconductor substrate, and the router optimizes the routing of input/output lines and the connection between standard cells so that the integrated circuit layout does not become overly congested by input/output and other routing lines. Other CAD tools are also available and may be used in the design of standard cell semiconductor IC devices.\nIn addition to the standard cells that may be selected from a library, the device design may include other standard cells such as user specified target cells with different operational characteristics. When such target cells are used along with the standard cells from the library to form a standard cell semiconductor IC device, additional challenges in the placement and routing of the cells is imposed. Target cells that operate at higher power consumption levels, for example, require the routing wires used to couple the target cells to ground and power sources, to operate at high power levels and carrying high currents. As such, these wires are especially susceptible to electromigration (EM) failure and even complete blowout of the wire, thus representing a significant shortcoming of conventional technology."}
{"text": "A current transformer monitors a current that passes through an electrical apparatus such as a power transformer, a generator, or a circuit breaker. To ensure that a current transformer is accurately monitoring the current, the metering or the relaying accuracy of the current transformer is verified. The metering accuracy determines the ratio error and the phase angle of the current transformer through a stringent certification process established by the Institute of Electrical and Electronics Engineers (IEEE) under IEEE C57.13 (1993). In comparison, the relaying accuracy relates to measuring the secondary excitation and resistance associated with the current transformer. Generally, current transformers intended for relaying applications do not undergo the certification process under IEEE C57.13 unless the current transformer is retrofitted for a metering application.\nBefore an accuracy test can be performed under IEEE C57.13 (1993), the current transformer is disconnected and removed from its position within the electrical apparatus. In a traditional accuracy test, a current, referred to as a primary current, circulates through the primary winding of the current transformer and a known standard transformer of the same current ratio. The ratio and phase displacement errors of the current transformer are then determined by comparison to the known standard. This process consumes a substantial amount of time and incurs a significant cost due to removing the current transformer from its installation on or within the electrical apparatus. Frequently, it is more cost-effective to replace rather than test the existing current transformer. It is therefore desirable to have a test method that overcomes the disadvantages associated with conventional methods for determining the accuracy of a current transformer."}
{"text": "The present invention relates to scanners, and more specifically to a carriage positioning structure for a scanner which enables the scanner to pick up the image of an object placed thereon, or turned upside down to pick up the image of an object placed below it.\nRegular scanners are commonly designed to pick up the image of flat sheets of document (2-dimensional scanning). If a sufficient depth of field is provided, a scanner can be used to pick up the image of a 3-D object. However, because regular scanners are commonly designed for 2-dimensional scanning, it is difficult to use a regular scanner to pick up the image of a 3-D object. When using a regular scanner to pick up the image of a 3-D object, the scanner may have to be turned upside down, however the image pick-up module (mechanism) cannot be stably reciprocated when the scanner is turned upside down."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit and an image processing apparatus having the same.\n2. Description of the Related Art\nAccording to a conventional technology for data transmission, a plurality of devices can be connected to a peripheral component interconnect (PCI) bus. One of the devices obtains a right to use the PCI bus, and such a device is determined as a master. The master device can transmit and receive data to and from any of the other devices.\nBecause the PCI bus employs a parallel transmission system, when the PCI bus receives a large amount of data, it is likely to cause a bottleneck. In order to overcome such a problem of the PCI bus, a technology called PCI Express (PCIe) has been proposed.\nThe PCIe allows point-to-point links between devices, and employs a serial transmission system.\nThe PCIe employs a full-duplex system in which transmission and reception of data are performed separately, and thereby the PCIe supports split transaction. Thus, the transmission and the reception of data can be performed at the same time. In the split transaction, a request and a response to the request are separately performed, and therefore a subsequent request can be issued before the response is obtained. In this manner, the transmission and the reception of data can be performed in an effective manner and at a high speed.\nIn the PCIe, a latency time of a read request to read data from a memory becomes longer. The latency time is a period starting from issuing the read request and ending at receiving the data from the memory. For this reason, a timing of issuing a read request affects the throughput efficiency of a bus in use for split transaction.\nFIG. 10 is a schematic diagram for explaining transmission efficiency of a conventional PCIe. When write requests are continuously received from an arbiter in an application specific integrated circuit (ASIC) or the like, the write requests are first issued, and afterwards a read request is issued. Therefore, a reception line is in an idle state during a period T1, during which write requests and data corresponding to the write requests are sent via a transmission line. As a result, the conventional PCIe obtains low transmission efficiency.\nFIG. 11 is a schematic diagram for explaining another example of transmission efficiency of a conventional PCIe. If the data transmission rate in an ASIC 5 is lower than the packet transmission speed of a PCIe 6, a request is sent to the PCIe after all of data corresponding to the request is stored in a data buffer that is arranged in a PCIe interface (I/F) circuit. Thus, it is possible to prevent data underflow. However, because a request cannot be sent to the PCIe before all the data is stored in the data buffer, the transmission line is in an idle state during a period T2. Occurrence of the period T2 causes the period T1 in the reception line to increase, and therefore the transmission efficiency becomes much lower.\nThe use of the PCIe makes it possible to perform the transmission and the reception of data separately and at the same time, nevertheless read data transmission can still be delayed.\nA master (bus master device) often executes some operation on data that is read from a memory, and then writes the resultant data in the memory. Therefore, a delay in time for receiving read data from the memory causes a delay in an issue of a write request by the master. As a result, the throughput of the whole system is lowered.\nAccording to a conventional technology disclosed in Japanese Patent Application Laid-Open No. 2006-154910, a memory controller selects a memory access request that allows continuous access to a memory upon sending a write request and a read request, and such continuous access achieves high-speed access to the memory. Furthermore, a plurality of read accesses is collectively processed as continuous read access, and a plurality of write accesses as continuous write access. A read cycle and a write cycle are alternately processed. Thus, it is possible to reduce the number of times of switching between a read access and a write access, and achieve high-speed access to the memory.\nHowever, such a conventional technology has a problem if employed in the PCIe. Specifically, when write requests are collectively issued via a transmission line, a reception line is in an idle state. On the other hand, when read requests are collectively issued, a transmission line is in an idle state. As a result, utilization efficiencies of the transmission line and the reception line are decreased, and the throughput of the whole system is lowered."}
{"text": "1. Field of the Invention\nThe present invention relates to a smoking article comprising a carbon monoxide reducing agent and a method of manufacturing the same, and a method of manufacturing the carbon monoxide reducing agent.\n2. Description of the Related Art\nCarbon monoxide that is generated during combustion of a smoking article is a substance that is preferable to be removed. Conventionally, attempts to reduce carbon monoxide included in mainstream smoke that is generated during smoking or in sidestream smoke during natural combustion have been done by incorporating a catalyst, an oxidant, an adsorbent and the like into a tobacco filler, a tobacco wrapping paper, a filter and the like.\nFor example, Jpn. PCT National Publication No. 2008-502343 discloses that silver oxide is used as a catalyst for oxidizing carbon monoxide. Jpn. PCT National Publication No. 2007-527782 and Jpn. PCT National Publication No. 2007-527698 disclose a catalyst comprising metal particles and/or a metal oxide of a nanoscale carried on a fibrous carrier or a high-surface area carrier. Jpn. PCT National Publication No. 2005-527205 discloses an oxyhydroxide compound that acts as an oxidizing agent for converting carbon monoxide to carbon dioxide, and exemplifies an oxyhydroxide compound of a transition metal and the like.\nHowever, use of a large amount of a noble metal or a transition metal as a metal causes increase in the manufacture cost of a catalyst. Furthermore, use of nanoscale particles may cause problems of difficult handling in the manufacture of smoking articles, increase in costs for pulverization, and the like. On the other hand, milliscale particles having a large particle size are difficult to be included into or carried by a smoking article in the manufacture.\nMeanwhile, Jpn. Pat. Appln. KOKAI Publication No. 2006-122793 discloses a catalyst including a composite oxide comprising calcium and aluminum as main components and also comprising iron as a catalyst for a shift reaction of water gas. In that publication, a catalyst comprising a composite oxide comprising calcium and aluminum (calcium aluminate) as a main component is manufactured by manufacturing a layered double hydroxide including calcium and aluminum in an alkali aqueous solution in advance, and sintering the product by heating.\nIt is known that calcium aluminate generally causes a hydration reaction when contacting with water and causes aggregation and crystal breakup. Since the catalyst in Jpn. Pat. Appln. KOKAI Publication No. 2006-122793 has a large BET specific surface area of 20 m2/g or more, it has a very large contact surface area with water, and thus a hydration reaction proceeds rapidly from the surfaces of the catalyst particles and causes breakup of the crystal structure or coagulation of the composite oxide. As a result, the catalyst particles are deactivated within a short time period."}
{"text": "As populations grow and more and more vehicles flood the roadways, parking locations have become problematic in terms of scarcity and proximity to desired sites, particularly in urban settings. Likewise, the presence and operation of large amounts of fuel-burning vehicles leads to undesirable increased air pollution levels, particularly in crowded urban settings. Additionally, difficulties in locating parking spaces in such locations generally requires drivers to undertake prolonged searching for such a purpose, thus resulting, with cumulative numbers of such vehicles facing such a necessary experience, in higher pollution levels on average.\nAs well, the capability of municipalities to afford proper infrastructure for vehicle parking is typically undertaken through metering devices, although such systems require constant monitoring to ensure proper payments have been made to the necessary effect. In tandem, then, there is a growing issue pertaining to the difficulties inherent in the availability and proper revenue generation of sufficient parking places in such urban locations. Parking garages, decks, commuter lots and other secluded parking lots have been constructed in the past to meet a certain degree of drivers' parking demands. However, these structures and areas are generally limited in terms of locations, typically cost a flat fee for entry, regardless of time spent, or, alternatively, require hourly payments of certain high levels that must be paid before egress is permitted, all in a manner to recoup the costs involved with the overall construction, necessitate either human or mechanical intervention to guarantee payments are made, and other machinery to prevent driver egress until payments have been received (such as gate arms, for example). In essence, although such structures and parking areas have proven successful to a certain degree, there remains a need for other alternatives to drivers who do not wish to be beholden to such strictures.\nIn particular, as alluded to above, street parking, rather than lot parking, has been an attractive alternative to many drivers, specifically as a means to allow for parking in a vicinity closer to a specific target location, reducing individual costs, and increasing freedom of movement through simple ingress and egress from a parallel parking space onto a roadway (as opposed to ingress into and egress from a parking facility that is controlled and requires some type payment prior to leaving). The greatest difficulty such street parking alternatives pose is the typical requirement for a driver to make numerous passes along certain streets to search for a any vacant space, let alone one that is of proper size to allow for such a driver to undertake a parallel parking operation without contacting adjacent vehicles, as well as to allow for sufficient space to facilitate subsequent egress from the target space. As an alternative to controlled parking facilities and areas, then, although street parking may be desirable to drivers, the inherent problems with such an unreliable possibility has made such a parking alternative not only difficult to manage, but also potentially dangerous as drivers take too many chances due to time constraints and/or impatience.\nThere are actually two different systems provided by municipalities with regard to street parking arrangements: marked spaces and unmarked spaces. Marked configurations provide separated and properly delineated locations for drivers to choose in comparison with unmarked types, for obvious reasons. Such marked spaces inherently provide not only set distances, locations, and metering possibilities for such municipal entities, but also, to the contrary, inefficient overall space usage and limited revenue generation as a result. San Francisco has begun a pilot program utilizing a sensor system to provide notifications to drivers of empty marked spaces as well as a relatively straightforward manner of implementing charging methods to generate revenue directly in relation to marked space usage. Such an overall system, however, is relatively straightforward itself in that notification of such available marked spaces is much easier through specific sensor utilization. As well, correlating such sensor technology to a metering system that permits indications of parking initiation and ultimately parking space evacuation merely requires indications of sensor operation in the limited areas of such marked locations. In other words, the implementation of a notification/metering system limited to specific delineated parking spaces involves mere sensor activation and deactivation specifically for the sensors present within the area related directly to the marked spaces themselves. Unfortunately, as alluded to above, such marked space usage is rather inefficient, particularly in terms of allowing for various vehicle sizes (i.e., compact vehicles may fit easily within a space whereas a moving van will require far more space, and thus, in certain municipalities may not be permitted to park on a street). The capability of avoiding such inefficiencies with available on-street parking spaces while also according a reliable method of notifying drivers of vacant space of suitable size for specific types of vehicles would not only allow for reduction of necessary paint or other indication materials for marked spaces, but also would permit more efficient space usage and, overall, a reduction in pollution through allowing drivers a more direct route to a possible parking space. Additionally, the capability of implementing a suitable metering and charging method in relation to such an unmarked parking space system would allow, ultimately, to effective and more efficient revenue generation as the amount of available parking spaces would be increased through the omission of demarcations. To date, however, such has not been accorded by the on-street parking industry.\nControlled parking environments (such as lots and garages at public events, airports, stadiums, commuter lots, office buildings or other large parking areas) have been developed with certain marked space notifications in the past, although even these methods have proven problematic and time-consuming to a certain degree due to the difficulties in identifying available parking space locations, particularly in lots that are partially filled or almost filled to capacity. Such parking facilities, however, have, at least in some circumstances, been outfitted with either manual counting means or sensor devices to provide indications to drivers in various ways as to vacant spaces, both in terms of number and, at times, locations. Such notification operations, however, are limited in their scope to such location and number issues; since each parking space present within the confines of outfitted parking facilities are pre-sized and provided specifically for vehicles of various sizes, and, as such, the facilities themselves are configured to compensate for specific numbers of spaces properly aligned for such vehicles to park appropriately in side-by-side relation (and to allow ingress and egress from such spaces easily as result), such monitoring and consequent notification of availability and location is relatively simple and straightforward. Likewise, there is no need nor concern with providing any notification basis with regard to parking availability for vehicles that are parked end-to-end instead in such situations as the overall lot design generally includes such side-by-side arrangements (with limited possibilities of different space sizes based on the potential for smaller or larger vehicles that may be present). In any event, these prior parking lot notification schemes have basically pertained to set parking designations with specific parking configurations that are relatively easy to monitor as to vacancies and locations as a result. To the contrary, street parking, particularly such parking omitting any space demarcations, is highly erratic in principle since end-to-end, parallel parking requirements depend almost entirely on varied vehicle lengths and proper placement of such vehicles without any set configurations in place. Marked on-street parking spaces at least provide specific locations for drivers to select; again, though, such marked spaces are rather inefficient through limiting the amount of available spaces to too great a degree. For instance, with a standard block in New York City is ⅛th of a mile, or 660 feet in length. With set spaces of 23 feet (with one foot between each marked space), roughly 27 standard vehicles may be parked (theoretically, as driveways, hydrants, and set distances from intersections would invariably limit such numbers to a smaller number, in actuality). The utilization of unmarked systems would thus increase the amount of available space by at least six (20 feet per vehicle instead of set 24 feet distances). Additionally, and as noted above, set limits of 23 feet (with another foot set between each space) would consequently limit the type of vehicle that may actual park within such set spaces. Unlined spaces would thus allow, again, for all types and sizes of vehicles to utilize such on-street possibilities. As such, a greater efficiency for on-street parking would be realized if marked spaces were not implemented for these reasons.\nEven with unrestrained on-street parking being the norm, other considerations may impact the overall capability of such a system. For instance, numerous items that prevent certain parking locations (such as fire hydrants, driveways, and the like) exist that may cause difficulties for a driver to properly find and select a parking spot. Additionally, the chaotic state of filling up limited spaces in end-to-end arrangements, particularly with, again, myriad lengths of vehicles (from small sub-compact vehicles to large trucks, vans, and/or buses, for example) makes predictability for a driver as to the proper space size for his or her vehicle rather daunting, beyond the initial consideration of whether any space is available to being with. In other words, in comparison with set parking lot configurations, the availability of spaces present on city streets is, suffice it to say, a difficult issue for a driver to overcome, particularly in terms of finding a suitable (and legal) space in a desirable location. In those situations, the longer a driver continues to seek for an available on-street parking space, the more fuel is consumed, the greater pollution is generated, and the more time lost for the driver in terms of his or her own efficiencies. With large numbers of vehicles searching for open spaces on a daily basis, the greater the accumulative effect such fuel consumption, undesirable exhaust, and lost time issues create.\nAdditionally, in terms of street parking issues, notifications are not easy to provide beyond the fact that end-to-end parking presents an extremely difficult prediction capability (due to the fact that vehicles do not have static lengths, not to mention the fact that other obstacles, such as fire hydrants, etc., may be present, as noted above) and the manner of providing such notifications, if they are provided in that context, would not be provided through, for example, typical board displays or other types of updated signage in a set location (such as is generally employed within parking facilities). The breadth of the street parking area in question would be too great to permit such a display to have any effect, basically, as the entirety of the area at issue would require street-by-street notifications to drivers in order to provide up-to-date and reliable information. As it is, a single, or perhaps two, three, or four, signs strategically placed in the multi-street area in question would require a driver to pass by, read, and comprehend the information provided in order to allow for such an on-street notification system to be effective. In essence, although street by street notifications may be effective in on-street parking configurations, the same has not been found true for typical signage utilized for parking lot situations.\nAs well, street metering requires investment in technology that requires expenditures for purchase, implementation, upkeep, and constant monitoring to be effective. As a revenue stream, such a system is rather straightforward (as discussed above) for set parking spaces (and thus set sensors with specific activation and deactivation sequences in relation to the presence of a vehicle in the target space). For unmarked spaces, to the contrary, such a needed overall system is, for lack of a better word, nonexistent. Without set spaces, it is difficult to have metering devices devoted to specific spaces (not to mention, individual meters are rather expensive to implement, operate, and monitor). Even with meters with marked spaces, such systems are rather cumbersome and unreliable (drivers do not always pay into the device, and thus an attendant must constantly monitor such areas for proper utilization; furthermore, limited time periods available to a driver for coin or other currency input would require repeated visits to the parking location to avoid incurring parking violations) and, as discussed at great length, inefficient in terms of available spaces and thus overall available numbers of vehicles that may be present within such spaces at any one time; optimizing and maximizing space availability and correlating such to revenue generating capacity would be far more attractive, for obvious reasons. Unfortunately, the capability of effectively charging drivers for unmarked street parking is limited to honor systems (for instance, requiring drivers to pay into permanent structures located on the street itself, and receiving a dashboard ticket for display during parking) as well as constant monitoring to ensure payment has been made for such a benefit. Avoidance of such cumbersome and relatively expensive requirements, in addition to the ability to provide effective and reliable notification to drivers of available and suitably sized spaces within an unmarked on-street parking configuration, has simply not been a possibility.\nOf additional importance is the ability to accord drivers of large trucks (18-wheelers, for example) reliable indications as to the availability of locations to rest during a long trip. Federal regulations require such drivers to keep specific logs as to the length of delivery trips as well as the actual time spent driving on a specific day in order to best ensure such driving activities are as safe as possible. To that end, truck drivers utilize any available location to pull over and rest if and when such is necessary during a long trip. Thus, as has been recognizable on interstate highways, it is common for multiple trucks to park end-to-end at highway entry ramps and exit ramps, sometimes in rather long lines. Even, however, if such parking configurations are commonplace, truck drivers have been known to park in isolation in such locations in order to get the necessary rest for a long journey, placing themselves in potential jeopardy of being hit by vehicles or even of invasions by criminals while they are asleep. Certain refilling stations offer amenities to truck drivers (showers, sleeping berths, etc.), as well, although such activities cost money to the user, and prior notifications of availability.\nInterstate (and other roadway) rest areas include segregated areas, generally, for such large vehicles to park, ostensibly to permit sufficient space dedicated for such a purpose. Within such highway arrangements, such rest areas are typically spaced from 50-80 miles apart, on average, thereby according such drivers the potential to properly time when they choose to stop. Unfortunately, however, the spaces available in such rest areas are limited and other than word-of-mouth from other drivers, there is little notification provided in terms of such parking space availability. As a result, this lack of notification (or parking availability guarantees) leaves such drivers at the mercy of chance in such situations; when they arrive, only then they can assess the parking space potential for a lengthy rest at a specific site. Additionally, however, it is not unheard of that smaller trucks (such as moving vans, for instance) may be parked in 18-wheeler designed (and typically dedicated) spaces, even though such vehicles may be better suited for the car/van locations at such rest areas. Again, this potential for loss of proper space for a large truck creates other scenarios that may exacerbate an already stressful driving experience.\nThus, there exists a need to provide not only on-street parking notifications as noted above, but also real-time parking space availability notifications for such large 18-wheel vehicles within rest area or other like locations with dedicated sites for such purposes. Furthermore, the ability to not only provide such availability indications but to monitor for vehicle sizes (to ensure proper utilization of such spaces for such specifically sized vehicles) as well as means to monitor the actual time spent within such spaces would be of enormous benefit for a variety of reasons within the driving industry. Additionally, the ability to provide a reservation service for drivers at specific locations would be extremely helpful, particularly to permit reliable timing for drivers to assess and undertake their required rest periods. To date, however, such a system has yet to be provided."}
{"text": "1. Field of the Invention\nThe present invention relates to methods for preparing 4-aminodiphenylamine (4-ADPA) and, more particularly, relates to a method for preparing 4-ADPA wherein aniline is reacted with nitrobenzene in the presence of a base, and under conditions wherein the amount of protic material, e.g., water, is controlled, to produce a mixture rich in the salt of 4-nitrodiphenylamine and/or the salt of 4-nitrosodiphenylamine. The 4-nitrodiphenylamine and/or 4-nitrosodiphenylamine salts can be isolated and subsequently hydrogenated or, alternatively, the reaction mixture itself is hydrogenated, to produce 4-ADPA in high yield. The present invention also relates to methods for preparing 4-ADPA intermediates as well as to alkylated p-phenylenediamine products useful as antioxidants.\n2. Related Art\nIt is known to prepare 4-ADPA by way of a nucleophilic aromatic substitution mechanism, wherein an aniline derivative replaces halide. This method involves preparation of a 4-ADPA intermediate, namely 4-nitrodiphenylamine (4-NDPA) followed by reduction of the nitro moiety. The 4-NDPA is prepared by reacting p-chloronitrobenzene with an aniline derivative, such as formanilide or an alkali metal salt thereof, in the presence of an acid acceptor or neutralizing agent, such as potassium carbonate, and, optionally, utilizing a catalyst. See, for example, U.S. Pat. Nos. 4,187,248; 4,683,332; 4,155,936; 4,670,595; 4,122,118; 4,614,817; 4,209,463; 4,196,146; 4,187,249; 4,140,716. This method is disadvantageous in that the halide that is displaced is corrosive to the reactors and appears in the waste stream and must therefore be disposed of at considerable expense. Furthermore, use of an aniline derivative such as formanilide, and use of p-chloro-nitrobenzene, requires additional manufacturing equipment and capabilities to produce such starting materials from aniline and nitrobenzene, respectively.\nIt is also known to prepare 4-ADPA from the head-to-tail coupling of aniline. See, for example, G.B. 1,440,767 and U.S. Pat. No. 4,760,186. This method is disadvantageous in that the yield of 4-ADPA is not acceptable for a commercial process. It is also known to decarboxylate a urethane to produce 4-NDPA. See U.S. Pat. No. 3,847,990. However, such method is not commercially practical in terms of cost and yield.\nIt is known to prepare 4-ADPA by hydrogenating p-nitrosodiphenylhydroxylamine which can be prepared by catalytic dimerization of nitrosobenzene utilizing, as a reducing agent, aliphatic compounds, benzene, naphthalene or ethylenically unsaturated compounds. See, for example, U.S. Pat. Nos. 4,178,315 and 4,404,401. It is also known to prepare p-nitrosodiphenylamine from diphenylamine and an alkyl nitrate in the presence of excess hydrogen chloride. See, for example, U.S. Pat. Nos. 4,518,803 and 4,479,008.\nIt is also known to produce 4-nitrosodiphenylamine by reacting acetanilide and nitrobenzene in DMSO in the presence of sodium hydroxide and potassium carbonate at 80.degree. C. for 5 hours. See Ayyangar et al., Tetrahedron Letters, Vol. 31, No. 22, pp. 3217-3220 (1990). See also, Wohl, Chemische Berichte, 36, p. 4135 (1903) and Chemische Berichte, 34, p. 2442 (1901). However, the yield of 4-nitrosodiphenylamine is low and is therefore not commercially practical. Furthermore, such method requires utilization of an aniline derivative, namely, acetanilide, and therefore increases the cost of the starting materials.\nThe process of the present invention does not include a halide source and therefore eliminates the expensive halide removal from the waste stream. In addition, the process of the present invention is much less expensive in terms of manufacturing costs, as well as the cost of raw materials, because instead of the derivatives of aniline and nitrobenzene, the present method utilizes aniline and nitrobenzene directly."}
{"text": "1. Field of the Invention\nThe present invention relates to an information processing device, an information processing method, and a program. More specifically, the present invention relates to an information processing device, an information processing method, and a program to control display of information depending on a mode in which the device is used, in relation to the information processing device combining an image display function and a printing function.\n2. Description of the Related Art\nThere has been widely used a digital photo frame as a photo frame type of a display device for displaying, for example, pictures taken with a digital camera. The digital photo frame, for example, is loaded with a recording medium such as a flash memory that stores an image taken with the digital camera, and displays the image stored in the medium on a display unit including a liquid crystal display (LCD) or the like. A user can use the digital photo frame as a photo frame by selecting a favorite picture from among a large volume of image data stored in the medium and displaying it on the display unit. Such a photo frame is described in, for example, Japanese Unexamined Patent Application Publication No. 2009-45139.\nThere is also a device configured to be connected to a printer via a universal serial bus (USB) terminal and a USB cable and output the image displayed on the display unit to the printer. Such a device with a printing function receives various settings for printing, command inputs, and the like through a remote control, a touch screen, or an operation switch or an operation button arranged in a space around the display unit.\nFor use as a photo frame, however, the switch or the button for printing presented on or around the display unit spoils a view disadvantageously causing disfigurement and lowering quality. Even if the operation button for printing can be provided on the back of the display unit, such a configuration is disadvantageous because it is difficult to operate the button when the user wants to print out the image while viewing the image displayed on the display unit."}
{"text": "1. Field of the Invention\nThis invention relates to methods and apparatus for processing and cleaning a substrate, and more specifically to methods and apparatus for cleaning semiconductor substrates after polishing of copper films.\n2. Background Information\nIn the manufacture of advanced semiconductor devices, copper (Cu) is beginning to replace aluminum (Al) as the material for metallization. Cu has become desirable due to its lower resistivity and significantly improved electromigration lifetime, when compared to Al.\nOne process for Cu metallization uses a dual damascene approach. As illustrated in FIG. 1a, a dielectric layer 110 is deposited above a substrate 100. Dielectric layer 110 may be made up of materials such as silicon dioxide. Vias and/or trenches 120 are then formed in the dielectric layer 110, as illustrated in FIG. 1b. Vias/trenches 120 may be formed, for example, using dry etching techniques. Next, a thin layer of barrier material (barrier layer) 130, for example, tantalum (Ta), titanium (Ti), or titanium nitride (TiN) is deposited as illustrated in FIG. 1c. After barrier layer 130 is deposited the vias/trenches 120 are filled with copper (Cu) layer 140, as illustrated in FIG. 1d. Copper layer 140 may be deposited using well known deposition techniques such as, for example, chemical vapor deposition (CVD), physical vapor deposition (PVD), or electroplating. In order to isolate the copper interconnects, as illustrated in FIG. 1e, the excess copper layer 140 and barrier layer 130 must be removed.\nOne method for removing the excess copper layer 140 and barrier layer 130 is polishing the surface of the substrate, for example, polishing using chemical mechanical polishing (CMP). In a CMP process, the semiconductor substrate is polished with a slurry containing abrasive particles, such as alumina particles, and an oxidant, such as hydrogen peroxide. In the CMP process, contaminants are introduced which include particles and/or metal contamination on the copper layer 150, dielectric surface 160, and in the dielectric subsurface 165.\nRegardless of how the CMP process is performed, the surface of semiconductor substrate must be cleaned of contaminants. If not removed, these contaminants may affect device performance characteristics and may cause device failure to occur at faster rates than usual. Cleaning the semiconductor substrate after chemical mechanical polishing of copper may be necessary to remove such contaminants from the copper layer and dielectric layers.\nOne method for cleaning the semiconductor substrate after polishing of the copper layer is brush scrubbing. Brush scrubbing, whether single-sided or double-sided brush scrubbing, is the industry standard for cleaning oxide and tungsten CMP applications. However, there are several problems associated with applying brush scrubbing to post copper CMP cleaning.\nOne such problem is brush loading. During the CMP process, the top surface of the copper layer may be oxidized and forms copper oxide, for example copper oxide (Cu.sub.2 O or CuO) or copper hydroxide (Cu(OH).sub.2). In basic or neutral pH cleaning environments, the copper oxide or copper hydroxide does not dissolve and may be transferred to the brushes, thus loading the brushes. The contaminated (or loaded) brushes may then transfer the copper oxide or copper hydroxide contaminants to subsequently processed substrates during cleaning.\nFor tungsten and other oxide applications, brush loading could be curtailed by adding a dilute ammonium hydroxide (NH.sub.4 OH). In the presence of NH.sub.4 OH, part of the copper oxide may form Cu(NH.sub.3).sup.2+ complex and may be dissolved; however, due to the high pH environment, the dilute ammonium hydroxide has been found to be insufficient to prevent brush loading of copper oxide. Additionally, it has been found that scrubbing with dilute ammonium hydroxide also causes etching of the copper layer and may cause serious surface roughening.\nBrush loading may also occur when alumina particles are used in the copper CMP process. In neutral or inorganic acid (e.g. HCl) cleaning environments, there is an electrostatic attraction between alumina particles and the silicon dioxide surface which makes it difficult to remove the alumina particles from the surface of the dielectric material. Because of the electrostatic attractive force, the alumina particles may also adhere to the brush and cause another brush loading problem with similar effects to those discussed above.\nYet another problem caused by the CMP process is that the surface and subsurface of the dielectric layer may become contaminated during polishing with metal from the copper layer and barrier layer as well as other contaminants from the slurry. During the CMP process, contaminants, especially metal contaminants, may penetrate into the dielectric layer up to approximately 100 angstroms (.ANG.) from the surface. Again, these contaminants may affect device performance characteristics and may cause device failure.\nThus, what is needed is a cleaning environment and methods for cleaning post copper CMP substrates that alleviate the problems of brush loading without affecting the quality of the copper and dielectric layers. Furthermore, what is needed is a cleaning environment and methods for cleaning post copper CMP substrates that have the capability of removing surface and subsurface contaminants from the copper and dielectric layers."}
{"text": "Field of the Invention\nThe present invention relates to wireless communications, and more particularly, to a method and apparatus for controlling an uplink power in a wireless communication system.\nRelated Art\nUniversal mobile telecommunications system (UMTS) is a 3rd generation (3G) asynchronous mobile communication system operating in wideband code division multiple access (WCDMA) based on European systems, global system for mobile communications (GSM) and general packet radio services (GPRS). A long-term evolution (LTE) of UMTS is under discussion by the 3rd generation partnership project (3GPP) that standardized UMTS.\nThe 3GPP LTE is a technology for enabling high-speed packet communications. Many schemes have been proposed for the LTE objective including those that aim to reduce user and provider costs, improve service quality, and expand and improve coverage and system capacity. The 3GPP LTE requires reduced cost per bit, increased service availability, flexible use of a frequency band, a simple structure, an open interface, and adequate power consumption of a terminal as an upper-level requirement.\nTo increase the capacity for the users' demand of services, increasing the bandwidth may be essential, a carrier aggregation (CA) technology or resource aggregation over intra-node carriers or inter-node carriers aiming at obtaining an effect, as if a logically wider band is used, by grouping a plurality of physically non-continuous bands in a frequency domain has been developed to effectively use fragmented small bands. Individual unit carriers grouped by carrier aggregation is known as a component carrier (CC). For inter-node resource aggregation, for each node, carrier group (CG) can be established here one CG can have multiple CCs. Each CC is defined by a single bandwidth and a center frequency.\nIn LTE Rel-12, a new study on small cell enhancement has started, where dual connectivity is supported. Dual connectivity is an operation where a given UE consumes radio resources provided by at least two different network points (master eNB (MeNB) and secondary eNB (SeNB)) connected with non-ideal backhaul while in RRC_CONNECTED. Furthermore, each eNB involved in dual connectivity for a UE may assume different roles. Those roles do not necessarily depend on the eNB's power class and can vary among UEs.\nUplink power control determines the average power over a single carrier frequency division multiple access (SC-FDMA) symbol in which the physical channel is transmitted. Uplink power control controls the transmit power of the different uplink physical channels. Efficient uplink power control method for CA or dual connectivity may be required."}
{"text": "Computing devices generally operate using a combination of firmware and software that is configured to execute on the hardware of a given computing device. A given set of firmware and software (e.g., the operating system) that is used to operate a computing device may be referred to as an operating image for the computing device.\nA computing device may come with a trusted operating image (e.g., an image that is signed and may be verified with a key) that was installed onto the computing device in the factory or updated using a trusted image source. However, computing device users may also wish to load and boot images that are not verified or not trusted (e.g., unsigned or self-signed images). For example, a user may wish to customize and make changes to the operating system for the computing device or the user may wish to install a different operating system in addition to or to replace the existing operating system."}
{"text": "1. Field of the Invention\nThe present invention relates in general to computing systems, and more particularly, to systems and methods for background destaging storage tracks from cache when one or more hosts coupled to the cache are idle.\n2. Description of the Related Art\nOne goal of computer storage systems is to reduce the number of destage conflicts when destaging storage tracks from a write cache so that the storage system operates more efficiently and/or faster. A destage conflict can occur when a storage track is being destaged from the write cache while a host is trying to write data to the storage track. This situation may be created because contemporary storage systems typically destage storage tracks shortly after the storage track is written to and the host must wait until the storage track is destaged from the write cache before the host is able to again write to the storage track. One technique to reduce the number of destage conflicts includes keeping storage tracks in the write cache for a longer period of time before destaging the storage tracks so that a storage track may be written to multiple times before the storage track is destaged. While this is an effective technique for reducing destage conflicts, it is also desirable that storage tracks not reside in the write cache too long so that the data in the storage tracks does not become stale."}
{"text": "Controlled sharing of objects is fundamental to distributed systems. However, only limited and simple means of sharing are supported in network or cloud-based systems. Current software systems are often constructed as a decentralized graph of collaborative, loosely-coupled services that include different protection domains, communication channels, execution environments and implementations.\nAs such, security and access control are of concern. Instead of using public-key based certificates or other sophisticated key agreement, cloud-based services typically use a pattern of pure bearer tokens or authentication credentials. Bearer tokens or credentials that grant unconditional authority are often sent directly between clients, servers and protection domains, even on unencrypted channels."}
{"text": "It is known for a rocket engine to include as a component thereof a thrust chamber having several sections joined to each other in an axial direction. Each of the sections is formed by a wall structure having an inner wall, an outer wall parallel thereto, and cooling channels formed between the inner and outer walls. The wall structure is continuous in the circumferential direction of the section.\nThe inner wall of each of two joined (to be joined) sections projects further (longer) in the extension direction of the wall structure than does the outer wall. The projecting end portion of the inner wall of one section is joined to the adjacent projecting end portion of the inner wall of the other section by a weld joint. In this way, a substantially continuous inner wall is achieved.\nThereafter, a ring-shaped element is radially arranged outside the weld joint, and said element is joined to the end portions of the adjacent outer walls. In this way, the cooling channels of one of the sections can communicate with the cooling channels of the adjacent section.\nEven though the above described rocket engine component performs satisfactorily, there is still a desire to increase the life of the component so that it can be used for an increased number of engine cycles."}
{"text": "Insect pests are a serious problem in agriculture. They destroy millions of acres of staple crops such as corn, soybeans, peas, and cotton. Yearly, these pests cause over $100 billion dollars in crop damage in the U.S. alone. In an ongoing seasonal battle, farmers must apply billions of gallons of synthetic pesticides to combat these pests. However, synthetic pesticides pose many problems. They are expensive, costing U.S. farmers almost $8 billion dollars per year. They force the emergence of insecticide-resistant pests, and they can harm the environment.\nOther approaches to pest control have been tried. In some cases, crop growers have introduced “natural predators” of the species sought to be controlled such as non-native insects, fungi, and bacteria like Bacillus thuringiensis. Alternatively, crop growers have introduced large colonies of sterile insect pests in the hope that mating between the sterilized insects and fecund wild insects would decrease the insect population. Unfortunately, success has been equivocal and the expense considerable. For example, as a practical matter, introduced species rarely remain on the treated land—spreading to other areas as an unintended consequence. Predator insects migrate, and fungi or bacteria wash off of plants into streams and rivers. Consequently, crop growers need more practical and effective solutions.\nOne relatively recent solution has been to genetically engineer crops to express plant lipases that have insecticidal properties. Until now, such insecticidal lipases have only been described in certain plants, such as patatin from the potato (U.S. Pat. No. 5,743,477) and pentin from the oil bean tree (U.S. Pat. Nos. 6,057,491 and 6,339,144). However, plant-derived lipases have the inherent disadvantage of having induced natural selection pressure in insects feeding on these plants in the wild. Thus, alternative lipases are needed for insect resistance management. The present invention is useful for avoiding the inherent disadvantage of pre-existing natural selection pressure, while conferring numerous other advantages such as low cost relative to repeated-application pesticides and effective insecticidal properties."}
{"text": "The following text should not be construed as an admission of knowledge in the prior art. Furthermore, citation or identification of any document in this application is not an admission that such document is available as prior art to the present invention, or that any reference forms a part of the common general knowledge in the art.\nTouchless automated hand-washing devices that incorporate rotating cylinders are known and have been in use for some time. These devices are designed to wash and/or provide a proper amount of anti-microbial solution to a person's hand and forearm within in a set time. The cylinders can be mounted in a free standing device and are adapted to receive the person's hand and forearm. The cylinders include an annulus or inner cavity that, in operation, is filled with cleaning fluid. As the inner cavity is filled, the cleaning fluid is forced out of a nozzle arrangement on the inner surface of the cylinder and into the interior of the cylinder. In order to provide greater coverage of the hand and forearm, the cylinders are rotated to provide a spray pattern.\nA nozzle arrangement known in the prior art is disclosed by U.S. Pat. No. 4,817,651 (“the '651 patent”). This patent discloses a three-set grouping of nozzles positioned on the interior surface of the rotating cylinder. A first set of nozzles is positioned on a closed end of the cylinder opposite from an open end through which a person's hands are placed. The closed end includes a frusto-conical shape that allows the first set of nozzles to be positioned such that cleaning fluid is sprayed toward the user's hand. Additionally, the frusto-conical end portion may include a drain through which spent cleaning fluid may exit the cylinder. A second set of nozzles are arranged in a ring on the inner surface of the cylinder just inside the open end. The second set of nozzles are angled downward, such that their spray is directed into the cylinder. A third set of nozzles are positioned in a helical array along the length of the cylinder.\nAs described in the '651 patent, the three-set nozzle arrangement has several advantages. Firstly, the fingernails and the ends of the fingers receive a direct spray from the first set of nozzles. Secondly, the spray pattern that results from the second set of nozzles produces a “curtain” that prevents the cleaning fluid from being splashed or sprayed out of the cylinder. Lastly, the helical arrangement of nozzles along the length of the cylinder in combination the rotation of the cylinders results in debris and spent cleaning fluid being swept downward, toward the closed end of the cylinder. In particular, the third set of nozzles are disposed in a left-hand helical pattern which, when combined with the clock-wise rotation of the cylinders, results in a succession of sprays that travel down the arm from above the wrist towards the fingertips.\nThe spray pattern that results from the nozzle arrangement depends in part on the orientation and spray pattern of the individual nozzles. While the '651 patent contains little discussion of these details, U.S. Pat. No. 5,823,447 (“the '447 patent”) is directed to an angled fan nozzle for use with a rotating cleaning cylinder. The nozzle construction disclosed in the '447 patent produces a flat “fan” pattern. In particular, cleaning fluid is ejected into the interior of the cylinder from a nozzle having a small square aperture. The spray of cleaning fluid spreads out through an angle while remaining substantially in a single plane. The nozzle is constructed such that the plane of the fan spray pattern is angled at 15 degrees with respect to the axis of the nozzle. As described by in the '447 patent, the nozzle may be positioned on the interior of the cylinder at a 15 degree tilt towards the closed end of the cylinder. The 15 degree angle of the fan spray combines with the 15 degrees tilt of the nozzle, resulting in a fan spray directed partially downwards, towards the closed end of the cylinder. With respect to the plane of the cylinder opening, the fan spray is directed downward at a 30 degree angle. This downwardly angled fan spray operates to direct spent cleaning fluid downwards, towards the closed end of the cylinder.\nWhile the '447 patent discloses an angle of the flat fan spray with respect to the plane of the cylinder opening, this reference is silent as to the rotational orientation of the fan spray pattern around with respect to the normal of the inner wall (or “roll” angle as defined later herein). For example, a particular nozzle may be positioned in a 0 degree rotational orientation, such that the flat fan spray is ejected from the nozzle “horizontally,” that is, in a plane parallel with the plane of the cylinder opening. Alternatively, a particular nozzle may be positioned in the 90 degree orientation such that the flat fan spray is ejected “vertically,” that is, in a plane perpendicular to the plane of the cylinder opening. As can be appreciated, a particular nozzle can be positioned at any angle between 0 and 90 degrees such that the flat fan spray has both a vertical and horizontal component.\nIt has been found that particular orientations and arrangements of nozzles not disclosed or suggested by the '651 patent or the '447 patent, such as the rotational orientation of particular nozzles, produce a novel spray pattern having increased coverage of the hand and forearm of the user. Additionally, the prior art fails to disclose other novel features associated with the cylinders that allow for faster and more reliable completion of wash cycles, safer operation, and easier assembly.\nNothing herein is to be construed as an admission that the present invention is not entitled to antedate a patent, publication or invention by another by virtue of prior invention."}
{"text": "The subject matter described herein relates generally to wind turbines and, more particularly, to blades for use on a wind turbine.\nMany known wind turbines include a tower and a rotor mounted on the tower using a nacelle. The rotor includes a number of blades that rotate to drive a generator through a gearbox using a rotor shaft, and the gearbox steps up the inherently low rotational speed of the rotor shaft such that the generator may convert the mechanical energy to electrical energy.\nThe swept area of the rotor is directly correlated with the power output of the wind turbine. Given that the swept area is essentially equal to the area of a circle (as calculated by the formula πr2) for which the length of each blade essentially defines the radius, an increase in the blade length tends to increase the power output potential of the wind turbine, and a decrease in the blade length tends to decrease the power output potential of the wind turbine. As such, increasing the blade length is sometimes desirable. However, the blades are often manufactured away from the installation site of the wind turbine, and blade transportation costs have been known to constrain the blade length increase and, therefore, the swept area and the power output potential of a wind turbine."}
{"text": "1. Field of the Invention\nThe present invention relates to methods of manufacture of neckties, and, particularly, means and methods for manufacturing an improved apparatus to anchor a four-in-hand necktie to a shirt so that it remains properly restrained and centered when adorned.\n2. Description of the Prior Art\nA conventional \"four-in-hand\" necktie is tied about the neck with a knot at a shirt collar and generally has an outward facing section passing down the shirt front to conceal shirt buttons and a somewhat narrower tail section which hangs down from the knot and is concealed by the outward facing section. Normally, a wearer might use one of a number of possible methods to control a tie, including: letting both the outward section and the tail section hang free; passing the tail section through a manufacturer's label regularly provided on the back of the outward facing section so that the two sections of the tie hang as a unit; or attaching the two sections to the shirt by use of various forms of jewelry, such as a tie-tack, a tie bar or a tie pin to keep the tie in place.\nThe problem with the non-restrained methods of wearing a tie mentioned above is that the tie sections hang free. This not only regularly presents disarray in appearance, but the tie sections frequently interfere with work, are soiled or ruined by dragging through food or other staining material, and may be quite hazardous if accidentally caught in a doorway or machinery.\nThe use of jewelry holding devices is common, but also suffers from a number of drawbacks. First, these devices often pass in and out of favor depending on preference and fashion trends. Second, they are often expensive and regularly cause damage to the fabric by creating permanent holes, creases, and blemishes in the tie. Finally, these jewelry devices anchor the tie to the shirt and often do not permit vertical movement of the tie relative to the shirt. Torso movements of the wearer cause the anchored tie to pull on the knot causing displacement of the knot and general disarray so as not to present an elegant appearance.\nNot surprisingly, a number of solutions have been proposed to attempt to avoid the above problems, but none is believed wholly satisfactory. U.S. Pat. No. 4,920,579 to Swain describes an apparatus to restrain necktie tails of a four-in-hand tie at the front of the shirt by relying on the manufacturer's sewn-in-place label on the back of the outward facing section and a flat plastic loop threaded through that label and attached to buttons on the shirt front. The plastic loop is somewhat stiff and has button holes at each end with a fixed distance between them. This arrangement is believed to have numerous drawbacks, including that it is awkward to employ--requiring keeping track of the loop, re-positioning the loop each time a new tie is worn, and often re-tying of the tie to center the device properly. Further, this device relies on the distance between buttons fixed by the manufacturer of the shirt which is not always the same from shirt to shirt. Moreover, this device may cause tie damage through constant pulling on the manufacturer's label--a use for which the label is not designed. Finally, the stiff plastic loop may prove uncomfortable to wear.\nAnother approach is described in U.S. Pat. No. 4,827,576 to Prince, Jr. There a button-slot neck tie fastener is permanently affixed to the manufacturer's label and a stiff portion having button slots is affixed to a shirt button. Apparently the tie's vertical movement relative to the shirt is permitted by the slot engaging the button. This device is believed to suffer from most of the same drawbacks as the Swain device, and may have the additional problem of possibly wearing out the shirt front where it contacts the device and the button stitching--thus possibly damaging both the tie and the shirt.\nStill another solution to keeping a conventional tie in proper position is disclosed by U.S. Pat. No. 4,972,523 to Begg. Begg contemplates the use of complementary VELCRO.RTM.-type loop-and-hook material, with a first element thereof on the back of the outward facing section and a second element thereof on a loop around the tail section. The second element has a slot on its opposite side for engagement with a shirt button. The two hook-and-loop elements engage each other to hold the tie in position and the tail section of the tie is threaded through the loop. Although this arrangement solves some of the problems presented by the Swain and Prince devices, its separate elements are believed to be somewhat awkward to use and the loop element may be prone to loss. More importantly, this design does not allow vertical movement of the tie relative to the shirt to follow movement of the torso of the wearer--leading to some of the same presentation problems as conventional jewelry restraining devices.\nApplicant's device disclosed in the parent application solves all of these problems. The invention disclosed in the parent application employs a tie restraint apparatus with a vertical member attached at its top and bottom to the back of the front section of the tie, and a horizontal member which surrounds the vertical member and slides between the vertical member's top and bottom positions. The horizontal member includes one or more openings in it which attach to buttons on the shirt.\nIn operation, once the necktie is tied in a conventional manner, the horizontal member is easily slid into a position on the vertical member which aligns with an opening with a button on the shirt. The button is then secured within the opening to provide an anchor for the tie. The concealed tail section of the tie may be inserted through the horizontal member to keep it safely secure. This entire apparatus may be permanently affixed to a tie and constructed from material identical or similar to the tie's so that it is completely concealed from vision and its presence is not evident to the wearer.\nThe attachment of the horizontal member to the shirt has proven to be a very secure method of retaining the tie in place. This is due in part to the free movement of the horizontal member on the vertical member which assures that the tie always stays in a straight and low-stress orientation. Moreover, the invention of the parent application is extremely simple and reliable in operation, may be instantly attached without complicated positioning of the tie relative to shirt buttons, may be inexpensively incorporated as original equipment on any conventional tie, and, due to its automatic adjustability, always presents an orderly appearance regardless of a wearer's torso movement.\nDespite the many advantages of the invention of the parent application, further improvements are considered to be possible, particularly in the method of manufacture. Although the invention of the parent application may be constructed from virtually any material, it has been determined that the apparatus functions significantly better if the material is stiffened to provide for smoother sliding between the vertical and horizontal members. Although this can be readily accomplished by employing a relatively stiff fabric for the apparatus, this tends to make the apparatus stand out from the tie in an undesirable manner.\nA problem which then emerges is how to employ the same material as the tie itself, so as to camouflage the apparatus, while providing a suitably stiff apparatus. Experimentation with lining the tie material has proven somewhat effective to this end, but is viewed to be far too cumbersome for full scale production.\nAccordingly, it is a primary object of the present invention to provide a fabric apparatus for maintaining a necktie in place which is stiffened to assure smooth operation.\nIt is another object of the present invention to provide a method of manufacture of such an apparatus which readily provides full camouflage of the apparatus during use.\nIt is a further object of the present invention to provide a method of manufacture of a necktie restraining means which is as simple and straightforward as possible to be readily adapted for full scale production.\nThese and other objects of the present invention will become evident from review of the following specification."}
{"text": "It is common in the processing of silicon semiconductor wafers to remove or etch portions of a layer of silicon dioxide using an etchant. A common etchant is hydrofluoric acid. This is typically done using an immersion process wherein a wafer, or one or more carriers containing a number of wafers, are immersed in the desired hydrofluoric acid containing processing fluid.\nOne of the disadvantages of immersion etching processes is that the wafers typically exhibit an increase in the numbers of particulates which become adhered to or imbedded in the wafer. As the feature size of semiconductor wafers continues to decrease into the sub-micron sizes, now typically in the range of 0.1-0.5 micron, the need for minimizing contamination becomes even more acute. These small feature sizes also create additional problems for immersion processing because surface tension effects can reduce etching uniformity and resulting product quality.\nIt has previously been recognized that silicon wafers can be processed by using etchant gases, including gases which contain hydrofluoric acid (HF). In an article entitled, \"Etching of Thin SiO.sub.2 Layers Using Wet HF Gas\" authors P. A. M. van der Heide et al., describe the etching of silicon dioxide layers by using vapor mixtures of HF, water and nitrogen. The article describes using a flow of nitrogen carrier gas through a vessel containing a 10% hydrofluoric aqueous solution. The nitrogen carrier produced a flow of etchant gas which was directed through a nozzle against the surface of a small silicon wafer being etched. A flow of dry nitrogen was passed over the opposite side of the wafer to reduce the effects of ambient atmospheric water which was present in the essentially open atmospheric process. Temperatures from 25.degree.-40.degree. C. were indicated for the HF solution and from 25.degree. C. to about 60.degree. C. for the wafer. The authors also reported that high temperature treatment under vacuum conditions following the etch provided complete removal of oxygen.\nOne deficiency of the process described is the continuation of etching after the intended process exposure period has finished. This results in uneven etching rates and problems in the resulting integrated circuit devices due to variations in feature sizes across the device. Etching uniformity is also deficient in that the process arrangement may not provide assurance of homogeneous presentation of the reactant gas to the wafer being processed. Variations in the uptake of the vapors by the incoming carrier gas stream can result in instantaneous variations in the etchant gas stream which can affect processing. This deficiency may not have been as significant in the reported processing because of the substantial diluting effect of the carrier gas which also resulted in very slow etching rates.\nAnother serious deficiency of the processing described by van der Heide is that the etching rates are very slow and accordingly not practical from a commercial processing standpoint where high volume production rates must be maintained. Therein the authors describe the need to produce an initial wetting time and then periods of 3.5-5 minutes were reported as acceptable for removing very thin layers of only 1.2-3.5 nanometers. The resulting indicated etching rates were about 1 minute per nanometer (1 minute per 10 Angstroms). These relatively slow etching rates would provide commercially unfeasible processing times of 5-100 minutes where layers of 50-1000 Angstroms are to be removed. Such processing times are sufficiently slow to prevent acceptance of the van der Heide vapor processing as a substitute to the faster processing times possible using the prevalent immersion processing.\nU.S. Pat. No. 4,749,440 to Blackwood et al. shows a gaseous process for removing films from substrates which utilizes an anhydrous HF supply. The process is described to involve passing a flow of dry nitrogen over the wafer and then introducing a flow of reactive gas which is preferably an anhydrous hydrogen halide gas such as anhydrous hydrogen fluoride gas. A flow of water vapor laden nitrogen is also passed over the wafer before the anhydrous HF gas flow is begun and this is continued until after the anhydrous HF gas flow has been stopped. This type of processing leads to nonhomogeneous mixing occurring during the brief processing period of approximately 5-30 seconds. During the initial period exclusively water laden carrier gas is first introduced. This subjects the surface of the wafer to a high gradient moisture increase over a very short period of time. This is immediately followed by anhydrous HF gas which is incapable of reaching any effectively uniform homogeneous or equilibrium condition with the water vapor because of the prior introduction of the water vapor and the quickly terminated introduction of the HF which creates a highly reactive combination varying from point to point across the wafer surface. The resulting highly reactive but nonhomogeneous etchant gas is capable of high etch rates. Unfortunately, the results using this process have proven to be highly variable with nonuniform etching on the same device being a common problem as well as nonuniformity from one device to the next. Acceptance of this processing system was initially enthusiastic but has been nearly abandoned by chip manufacturers because of the seriousness of the nonuniformity problem in etching rates.\nThe nonuniformity problem also results from variations in the amount of water present in the matrix of the material being processed which can have very significant effects on the effective localized HF-water concentrations on the surface of the chip during the high speed etching reaction which occurs during this type of processing. Wafers otherwise processed in a similar manner may exhibit highly differing etch rates merely because they have been allowed to sit for hours under ambient conditions thereby uptaking atmospheric moisture to a substantially greater degree than other wafers processed soon out of a furnace or other moisture eliminating processing. Such variations in moisture content of the wafers are typical and any special pre-atmospheric treatment necessarily increases processing time, processing logistics, or both.\nAnother more recent approach to vapor etching of wafers using HF is incorporated into the processing machine referred to as the EDGE 2000 from Advantage Production Technology. This system utilizes a specially configured processing chamber which vacuum treats the wafer prior to processing with the etchant. This approach attempts to remove residual moisture from the wafer to address the nonuniform etching rate problem discussed above. The relatively short duration vacuum processing cannot remove all moisture content variations. The wafer is exposed to a highly reactive HF-water gas stream which is directed at the wafer in a vertical orientation from one or both sides. The jet of incoming reactant gas impinges upon the surface or surfaces of the wafer at localized central areas and typically results in nonuniformities in mass transfer due to this localized impingement despite the vacuum processing directed toward water removal.\nThe inventor hereof has discovered that vapor processing of semiconductor wafers, substrates and similar units using the novel processing described herein can result in the formation of undesirable particulate in some types of processing. The particulate has been manifested in localized deposits, or alternatively, in relatively diffuse layers generally deposited over the entire treated surface of the wafer. In some cases the particulate will be sufficiently concentrated to form a visible haze. In other cases the particulate will only be observable using magnifying instrumentation.\nThe exact nature of the particulate or particulates which form on the surface of the wafer is not known with certainty. However, it is believed to be one or more fluorosilicates formed by competing side reactions to the removal of silicon oxides by hydrofluoric acid. One possible fluorosilicate which may be formed is H.sub.2 SiF.sub.6, hydrogen hexafluorosilicate. Others compounds or solid phase mixtures may also be present or the primary particulate being formed.\nThus there remains a strong need in the art for gaseous HF and other chemical processing which will provide uniform and repeatable results with highly effective etching rates while achieving low contamination and particle counts. Other objectives and advantages of the invention are also indicated herein."}
{"text": "A radio apparatus that performs communications using a random access method, such as a wireless LAN, is required to keep awaiting a signal, since it is not known when the signal is transmitted. In general, a period of time for awaiting a signal is longer than period of time for receiving a signal. In order to reduce power consumption of the whole radio apparatus, it is important to reduce power consumption during standby time.\nAs a method for reducing power consumption of a radio apparatus, for example, a method for improving a configuration of a synchronizing circuit in a digital demodulating unit is known (see for example, Milos Krstic, et. al., “Optimized low-power synchronizer design for the IEEE 802.11a standard,” ICASSP 2003 pp. 333-336.).\nHowever, the method described in the Non-Patent Document 1 has a problem such that power consumption of an analog demodulating unit in a receiving unit, cannot be reduced while power consumption of a digital demodulating unit in the radio apparatus can be reduced.\nThe analog demodulating unit consumes as much power during standby time as when receiving a signal. Therefore, it is desired to reduce power consumption when the analog demodulating unit is in standby state in order to attain lower power consumption of the radio apparatus."}
{"text": "The present invention generally relates to the deglycerolization of blood, and more particularly, to a system which controls the deglycerolization of blood by monitoring the segregation of erythrocytes by size.\nThe Armed Services Blood Program Office (ASBPO) has established a policy of maintaining pre-positioned stockpiles of frozen red blood cells, and utilizing these stockpiles in times of conflict for U.S. combat casualties. In order to implement this policy, glycerol is allowed to be absorbed by red blood cells, which then are frozen and stored. The glycerol prevents damage to the erythrocytes. Presently, the only method approved by the Food and Drug Administration (FDA) for processing thawed-frozen red blood cells uses an open, nonsterile wash system that is manually monitored and operated. This system generally requires about 11/2 to 2 hours to thaw and deglycerolize red blood cells from a cryogenic state. Because this system is not sterile, the FDA mandates that thawed-frozen red blood cells processed this way must be transfused within 24 hours or discarded. However, the time restrictions and requirement to discard the blood are not compatible with the logistics of the ASBPO policy. Therefore, a need exist for a sterile, automated method for monitoring and controlling the deglycerolization of thawed red blood cells in a more timely manner compared to the processing time of the standard method."}
{"text": "1. Field of the Invention\nThe present invention relates to a computer power supply, and particularly to a computer power supply mounting apparatus.\n2. Description of the Related Art\nMany types of electronic devices, such as computers, include a power supply which is mounted in a housing or chassis along with the other circuits, components, and peripheral devices. Since the power supply receives alternating current electrical power and converts the alternating current to direct current for operating the various electrical components, safety regulations require that the power supply be completely enclosed in a case to prevent access to potentially harmful high voltages. In the past, most power supplies were mounted to the chassis using several screws, making it difficult and time consuming to replace or remove the power supply for repair.\nIn U.S. Pat. No. 5,790,372, an improved power supply mounting apparatus is disclosed. A computer includes a chassis and a power supply mounted in the chassis. One or more hooks are provided on the chassis and the power supply, so that the power supply can be placed in the chassis and moved relative to the chassis to engage the hooks and secure the power supply in the chassis against movement in a first plane. A screw is provided for engaging the chassis and the power supply for securing the power supply against movement in a plane perpendicular to the first plane. A tab on a wall of the chassis guides the power supply.\nIn U.S. Pat. No. 6,270,046, another improved power supply mounting apparatus is disclosed. A computer system includes a chassis which is formed by assembling front, rear and bottom panels together. An apparatus for mounting a switching power supply to the computer system includes a bracket and a drive cage with a pair of hooks. The drive cage bridges between the front and rear panels. The bracket includes a hinge portion and a supporting portion adapted to support the switching power supply. The hinge portion of the bracket is slidably engaged with the hooks for fixing the switching power supply to the rear panel of the chassis with screws.\nTherefore, what is needed is an screwless apparatus for mounting a power supply, thereby to enable easy and convenient installation and removal of the power supply."}
{"text": "An X-ray CT (Computed Tomography) apparatus calculates an X-ray absorption coefficient at each point from the measured projection data that is obtained by imaging a subject from multi-directions, and obtains as a tomographic image of a subject, an X-ray absorption coefficient distribution image (hereinafter, referred to as “CT image”) made up of plural pixels. The CT image acquired through this apparatus allows a patient's disease condition to be diagnosed accurately and immediately in medical practice, and it is clinically effective. However, a constant amount of radiation exposure is unavoidable in order to acquire an image of high quality necessary for diagnosis by a doctor. On the other hand, if the radiation dose is lowered in order to implement low dose scanning, a noise ratio increases to the signals being detected, and a large amount of linear streak artifact and/or graininess noise may occur, causing erroneous diagnosis. Therefore, it is desired to reduce such streak artifact and noise in the low dose scanning, and to achieve compatibility between high-quality diagnosis and low dose scanning.\nIn view of the situation above, the Patent Document 1 describes that there is employed an iterative approximate reconstruction method to iteratively correct a CT image so as to satisfy the condition that the calculated projection data is equal to the measured projection data, and reduce the noise. In the iterative approximate reconstruction method, firstly, a reconstructing process is performed on the region for reconstruction (hereinafter, referred to as “FOV”) including the subject, and an image (FOV image) is obtained. Next, projection data is calculated from the reconstructed image (FOV image) (hereinafter, referred to as “forward projection calculation”), and iterates the operation for correcting the CT image, in comparison to the aforementioned measured projection data. This may enhance the precision of the CT image within the FOV. The iterative approximate reconstruction method requires the condition that the FOV incorporates the subject therein, and therefore, if the correction of the subject is unrelated to the diagnosis, it may cause increase of the calculation amount.\nAs a publicly known analytical reconstruction method, an extensive reconstruction method is suggested to perform reconstruction operation on a local region, so as to reduce the calculation amount. In view of this, it is conceivable to apply the extensive reconstruction technique to the iterative approximate reconstruction method, and reduce the calculation amount. However, in the iterative approximate reconstruction method, it is difficult to perform the reconstruction operation only on the local region out of the measured projection data, and therefore, it is necessary to extract from the measured projection data, projection data being included in the local region (hereinafter, referred to as “local measured projection data”).\nThe Non Patent Document 1 discloses a technique that applies the extensive reconstruction technique to the iterative approximate reconstruction method. This technique reconstructs a CT image on the FOV incorporating the subject with a table, a fixture, and the like (hereinafter, referred to as “large FOV image”). Next, an image of a background region other than the local region is extracted from the large FOV image, and thus extracted image of the background region is subjected to the forward projection calculation, thereby obtaining background projection data. Thus obtained background projection data is subtracted from the measured projection data, thereby extracting local measured projection data, and calculating the CT image of the local region (hereinafter, referred to as “small FOV image”). Accordingly, it is possible to apply the iterative approximate reconstruction method to the local measured projection data, and improve the precision of the small FOV image.\nOn the other hand, in the analytical reconstruction method, when data obtained at some projection angles fails to be measured on the periphery of the FOV, it is known that the precision of the CT image of the FOV is remarkably deteriorated (incomplete reconstruction). The Patent Document 2 advises to estimate the subject within the region where data is unmeasurable on the periphery of the FOV, thereby improving the precision of a CT value (X-ray absorption value) in the FOV image."}
{"text": "The present invention generally relates to data processing. The invention relates more specifically to a mechanism that determines the order of execution among a plurality of code modules.\nComputer networks have become ubiquitous in the home, office, and industrial environment. As computer networks have grown ever complex, automated mechanisms for organizing and managing the networks have emerged. These mechanisms are generally implemented in the form of one or more computer programs, and are generically known as network management systems or applications.\nFIG. 1 is a simplified diagram of a network 100 that is managed by a network management station 10. The network 100 comprises one or more network devices 102, such as switches, routers, bridges, gateways, and other devices. Each network device 102 is coupled to another network device 102, or to one or more end stations 120. Each end station 120 is a terminal node of the network 100 at which some type of work is carried out. For example, an end station 120 is a workstation, a printer, a server, or similar device.\nEach network device 102 executes a network-oriented operating system 110. An example of a network-oriented operating system is the Internetworking Operating System (IOS) commercially available from Cisco Systems, Inc. Each network device 102 also executes one or more applications 112 under control of the operating system 110. The operating system 110 supervises operation of the applications 112 and communicates over network connections 104 using an agree-upon network communication protocol, such as Simple Network Management Protocol (SNMP).\nEach device 102 stores information about its current configuration, and other information, in a Management Information Base (MIB) 114. Information in the MIB 114 is organized in one or more MIB variables. The network management station 10 can send xe2x80x9cfetchxe2x80x9d and xe2x80x9csetxe2x80x9d commands to the device 102 in order to retrieve or set values of MIB variables. Examples of MIB variables include sysObjectID and sysDescr.\nPreferably the network management station 10 is a general-purpose computer system of the type shown and described further herein in connection with FIG. 10. The network management station 10 executes one or more software components that carry out the functions shown in block diagram form in FIG. 1. For example, the network management station 10 executes a basic input/output system (BIOS) 20 that controls and governs interaction of upper logical layers of the software components with hardware of the network management station. An example of a suitable BIOS is the Phoenix ROM BIOS. The network management station 10 also executes an operating system 30 that supervises and controls operation of upper-level application programs. An example of a suitable operating system is the Microsoft Windows NT(copyright) operating system. The network management station 10 may also execute other operating systems that may not require a BIOS 20, such as UNIX-type operating systems, microkernel-based operating systems, etc.\nThe network management station 10 executes an asynchronous network interface (ANI) 50 under control of the operating system 30. The ANI 50 provides an interface to the network 100 and communicates with the network using SNMP or another agreed-upon protocol. The ANI 50 provides numerous low-level services and functions for use by higher-level applications.\nThe network management station 10 executes a network management system 40 that interacts with a database 60 containing information about the managed network 100. The network management system 40 is an example of a network management application. Using a network management application, a manager can monitor and control network components. For example, a network management application enables a manager to interrogate devices such as host computers, routers, switches, and bridges to determine their status, and to obtain statistics about the networks to which they attach. The network management application also enables a manager to control such devices by changing routes and configuring network interfaces. Examples network management applications are Cisco Works, CiscoWorks for Switched Internetworks (CWSI), and Cisco View, each of which is commercially available from Cisco Systems, Inc.\nThe ANI 50 and network management system 40 need not execute or reside on the same physical computer. They may execute on different machines.\nA requirement of many network management applications is the ability to rapidly adapt to new kinds of network devices, such as switches, routers, and other hardware and software. Often, it is necessary to support a new kind of device that is introduced into the network, released or shipped by the device vendor other than when the network management application is updated, released or shipped. Thus, it is necessary to support new devices and new services when the network management application is installed in the field or running in a managed network.\nAccordingly, the network management application must accept new requirements in the field. In particular, the execution environment of the network management application must be able to accommodate new processes and services to be executed. Existing processes must be adapted to new devices.\nIn past approaches, there has been no simple way to integrate a new process, service, or software component and its execution into the existing network management application. Generally, in past approaches, the existing network management application must be patched or updated.\nBased on the foregoing, there is a clear need in this field for a network management system that can accept field updates of new executable processes or services.\nThere is also a need for such a system that can integrate the new executable processes or services into the existing execution flow of the network management system.\nIn particular, there is a need for a system that can determine, at execution time, the order of execution of a plurality of modules that represent a process or service to be executed with respect to one or more network devices.\nThere is also a need for such a system that can resolve the run-time execution order of a plurality of executable components based upon interdependencies of the executable components.\nThere is also a need for such a system that can resolve the run-time execution order of executable components based upon requirements such as time of execution, type of function to be executed, and resources that are required to be available before a particular component can execute.\nThe foregoing needs, and other needs and objects that will become apparent in the following description, are fulfilled in the present invention, which comprises, in one aspect, a method of dynamically determining an execution order of a plurality of executable components, the method comprising the steps of associating, with a first component of the plurality of executable components, one or more preconditions to execution of the first component; storing a partial order of execution of the plurality of executable components based on the pre-conditions; and generating a final order of execution of the plurality of executable components based on the partial order.\nOne aspect of this feature is that the step of associating further comprises the steps of associating a time base with the first executable component. According to another feature, the step of associating further comprises the steps of associating a precondition value that represents a category of service to be executed before the first executable component is executed. In another feature, the step of associating a precondition value further comprises the steps of associating a time of execution and a category with the first executable component.\nAccording to another feature, the step of associating a precondition value further comprises the steps of associating, with a first service module function of an asynchronous network interface in a network management system, an execution time, a name of a second service module function, and a type of the second service module function. In another feature, the step of associating a precondition value further comprises the steps of associating a time of execution value indicating that the first executable component is to be executed during device discovery processing.\nAnother feature is that the step of associating a precondition value further comprises the steps of associating a time of execution value indicating that the first executable component is to be executed at a fixed pre-determined time. Still another feature is that the step of associating a precondition value further comprises the steps of associating a time of execution value indicating that the first executable component is to be executed upon request by an external process.\nIn another feature, the step of storing further comprises the steps of storing a directed acyclic graph in a memory, in which the directed acyclic graph comprises a plurality of nodes, each node comprising information representing one of the preconditions. Still another feature is that the step of generating a final order of execution further comprises the steps of conducting a breadth-first traversal to each of the plurality of nodes of the directed acyclic graph, and building the final order of execution based on an order in which the nodes are traversed.\nYet another feature is that each node is associated with a level of the graph, and further comprising the steps of deleting one of the nodes when such node is encountered at more than one level of the graph. Still another feature is that one of the pre-conditions is a unit value that represents a data set that is required to be available before the first executable component is executed."}
{"text": "The present invention relates to an axial-flow turbine, and more particularly to a turbine nozzle and a moving blade forming an fluid passage of the axial-flow turbine.\nA variety of techniques relating to the axial-flow turbine have been employed to improve an internal efficiency of the turbine so as to improve the performance of the same. Since a secondary flow loss among internal losses experienced with the axial-flow turbine is a loss of a type common to all stages of the turbine, a contrivance that is capable of preventing the above-mentioned loss has been required.\nFIG. 19 shows cross sections of turbine stages of a usual axial-flow turbine including nozzle blades and moving blades. Referring to FIG. 19, a plurality of nozzle blades 4 are radially secured between an outer diaphragm ring 2 and an inner diaphragm ring 3 which are fit to a turbine casing 1 so that a nozzle blade passage is formed. A plurality of moving blades 6 is disposed at a downstream side of the nozzle blade passage. Each of moving blades 6 is sequentially implanted in the outer surface of the rotor wheel 5 at predetermined intervals in the circumferential direction of the rotor wheel 5. The tip of the moving blades 6 are attached to a cover 7 so that leakage of working fluid is prevented. Both the nozzle blades 4 and the moving blades 6 form a working fluid passage of this stage of the turbine.\nA next (second) stage of the turbine, which is located at a downstream side of the above (first) stage, has a rapidly enlarged passage for the working fluid. This passage is composed of a nozzle blade passage and a moving blade passage as well as the above working fluid passage. The nozzle blade passage is formed by an outer diaphragm ring 8, an inner diaphragm ring 9 and nozzle blades 10. The moving blade passage is formed by both moving blades 12 implanted in a rotor wheel 11 and a cover 13 attached to the tip of the moving blades 12.\nIn the second stage, the working fluid expands from a high-pressure condition to a low-pressure condition trough the passage, so the specific capacity (volume) of the fluid enlarges. To correspond to such enlargement of specific capacity, the inner wall of the passage is inclined in such a manner that the area of the passage is enlarged in the downstream direction.\nThrough the above-mentioned passage of the two stages, the working fluid generate a secondary flow at the nozzle blades 4 and 10. This mechanism of generating the secondary flow will now be described with reference to FIG. 20.\nWhen the working fluid, which is high-pressure steam or the like, flows in the nozzle blade passage between the nozzle blades, the working fluid is curved into a circular arc form in the nozzle blade passage as indicated with a two-dot chain line shown in FIG. 20. At this time, centrifugal components are generated in a direction from an extrados E of the nozzle blade 4 to an intrados F of another nozzle blade 4. Since the centrifugal components and the pressure in the nozzle blade passage are in equilibrium, the static pressure at the intrados F of the nozzle blade 4 is raised.\nOn the other hand, the pressure at the extrados E of the nozzle blade 4 is lowered because the flow velocity of the working fluid is high along the extrados E. As a result, a pressure gradient is generated in a region of the nozzle blade passage from the intrados F of the nozzle blade 4 to the extrados E of the another nozzle blades 4. As shown in FIG. 20, also a pressure gradient of the foregoing type is generated between the inner wall of the root of the nozzle blades and a layer adjacent to the outer surface of the tip of the nozzle blades in which the flow velocity is low, that is, in the boundary layer. In the portions adjacent to the boundary layer, the flow velocity is low and the acting centrifugal component is weak. Therefore, the flow of the working fluid cannot withstand the pressure gradient generated in a direction from the intrados F of the nozzle blade 4 to the extrados E of the another blade. As a result, the flows are generated in a direction from the intrados F of the nozzle blade 4 to the extrados E of the another nozzle blade, as indicated with symbols f1 and f2 shown in FIG. 20. The flows f1 and f2 collide with the extrados E of the nozzle blade 4 and curl up. As a result, secondary flow eddies 14a and 14b are generated adjacent to the inner wall of the root of the nozzle blades 4 and the outer wall of the tip of the same.\nFIG. 21 is a diagram showing a mechanism of the moving blades 6 disposed downstream from the nozzle blades 4 to generate a secondary flow. Since the mechanism of the secondary flow the moving blades 6 is substantially the same as the mechanism of the nozzle blades 4 to generate eddies in the secondary flow, the same functions as those shown in FIG. 20 are given the same reference numerals and symbols. As can be understood from FIGS. 22 and 23 showing losses of the nozzle blades 4 and the moving blades 6, eddy losses are caused from the secondary flow eddies. Thus, excessive losses are produced in the portions adjacent to the inner and outer walls of the turbine blades.\nIf secondary flow eddies 14a and 14b are generated, a portion of energy of the working fluid is dispersed. Moreover, nonuniform flows of the working fluid are formed, thus causing a problem to arise in that losses of the nozzle blades and the moving blades are enlarged and the performance of the stages to deteriorate excessively.\nTo prevent the secondary flow loss caused from the secondary flow eddies 14a and 14b generated in the above-mentioned passage (stages), a variety of techniques have been researched and developed. For example, a nozzle blade having a reduced outer surface has been employed. This reduced outer surface has a irregularities formed in the tip of the nozzle blade to reduce the height of the flow passage in the downstream direction. FIG. 24 is a cross sectional view showing a turbine nozzle having the nozzle blade 15 having a reduced outer surface. The nozzle blade 15 having the reduced outer surface causes flows along the outer surface of the nozzle blade 15. Thus, the flow line is shifted toward the inside portion (toward the central portion) of the nozzle blade passage as indicated with an arrow ft. It leads to a fact that the flow lines in the central portion and the root (inside) portion, similarly to those along the outer surface, are shifted inwards (toward the central portion), as indicated by arrows fp and fr. As a result, the flow lines pushes the flows to the inner wall of the nozzle blade 15 in portions adjacent to the root of the nozzle blade 15. Thus, enlargement of the boundary along the inner wall can be prevented so that enlargement of the loss caused from the secondary flow eddies is prevented.\nFIG. 25 shows a distribution of losses reduced attributable to an effect of the conventional nozzle blade 15 having the reduced outer surface to prevent enlargement of losses caused from eddies in the secondary flow. As can be understood from FIG. 25, losses can significantly be reduced in the portions adjacent to the root of the nozzle blade. Improvement in the performance has been confirmed also in overall efficiency experiments of the turbine stages.\nThe fact that the nozzle blade 15 having the reduced outer surface is able to improve the performance has been confirmed in the above-mentioned stage efficiency experiments. However, local separation of the flow at the tip of the nozzle blade takes place attributable to a rapid shift of the flow line, as shown in the distribution of losses in the trailing edge of the nozzle blade. Therefore, the secondary flow cannot satisfactorily be improved.\nMoreover, a great portion of the working fluid flows adjacent to the root of the nozzle blade. Therefore, considerable change in the flow rate occurs in the direction of the height of the nozzle blade.\nTherefore, the stage performance realized by the nozzle blade 15 having the reduced outer surface can furthermore be improved. That is, a nozzle blade passage is required which is capable of preventing separation of flows of the working fluid and improving the flow rate characteristic at the tip of the nozzle blade.\nIn view of the foregoing an object of the present invention is to provide a turbine nozzle and a turbine moving blades of an axial-flow turbine capable of reducing a loss in the secondary flow with a simple structure.\nThis object can be achieved according to the present invention by providing an axial-flow turbine comprising: an outer diaphragm ring and an inner diaphragm ring forming together an annular fluid passage; and a plurality of nozzle blades disposed in the annular passage, each of the nozzle blades being formed into a warped shape such that a central portion in a lengthwise direction of the nozzle blade maximally projects in a downstream direction.\nIn preferred embodiments, the annular fluid passage has a stepped portion at an inner surface of the outer diaphragm ring and an outer surface of the inner diaphragm ring, the stepped portion having a curvature surface so that the height of the fluid passage is reduced in a downstream direction thereof.\nThe stepped portion has a height in a radial direction of the fluid passage, the height being described by the relationships:\n0xe2x89xa6h1/L1 less than 0.05\n0.1 less than h2/L1 less than 0.2\nwhere L1 is the height of an leading edge of the nozzle blades, h1 the height of the stepped portion provided for the inner diaphragm ring and h2 is the height of the stepped portion provided for the outer diaphragm ring.\nEach of said turbine blades has an axial distance from the leading edge of the diaphragm to the trailing edge of the nozzle blades, the distance being described by the relationships:\nZt less than Zr less than Zp\nwhere Zt is the distance at the outermost end of the nozzle blades, Zr is the distance at the innermost end of the same and Zp is the distance at the central portion of the same.\nThe height L2 of the nozzle blades at a trailing edge is made to be smaller than the height L1 of the nozzle blades at a leading edge (that is L1 greater than L2).\nThe fluid passage is structured such that the inner surface of the outer diaphragm ring and the outer surface of the inner diaphragm ring are outwards inclined in the downstream direction.\nAn angle of inclination of said fluid passage is descried by the relationships:\n0xc2x0xe2x89xa6xcex81 less than xcex83 less than xcex82\nwhere xcex81 is an angle of inclination of the outer surface of the inner diaphragm ring, xcex82 is an angle of inclination of the inner surface of the outer diaphragm ring in the leading edge of the nozzle blades and xcex83 is an angle of inclination of a portion of the inner surface of the outer diaphragm ring following the tailing edge of the nozzle blades.\nThe height L2 of the nozzle blades at an trailing edge is made to be larger than the height L1 of the nozzle blades at an leading edge (that is, L1xe2x89xa6L2).\nThe fluid passage is structured such that the inner surface of the outer diaphragm ring is outwards inclined in the downstream direction and the outer surface of the inner diaphragm ring is inwards inclined in the downstream direction.\nAn angle of inclination of said fluid passage is descried\nxcex81 less than 0xc2x0 less than xcex83 less than xcex82\nwhere xcex81 is an angle of inclination of the outer surface of the inner diaphragm ring is, xcex82 is an angle of inclination of the inner surface of the outer diaphragm ring in the leading edge of the nozzle blades and xcex83 is an angle of inclination of a portion of the inner surface of the outer diaphragm ring following the trailing edge of the nozzle blades.\nThe fluid passage is structured such that the cross sections of the nozzle blades at the tip and the root of the nozzle blades are shifted in the circumferential direction of the annular fluid passage.\nA throat width between adjacent two nozzle blades is determined by the relationships:\nSpxe2x89xa6Sr less than St\nwhere Sp is the throat width at the central portion in the lengthwise direction of the nozzle blades, Sr is that at the root and St is that at the tip.\nIn other side of the present invention, an axial-flow turbine comprising: an outer diaphragm ring and an inner diaphragm ring forming together an annular fluid passage; and a plurality of nozzle blades disposed in the annular passage, wherein said annular fluid passage has a stepped portion at an inner surface of the outer diaphragm ring and an outer surface of the inner diaphragm ring, the stepped portion having a curvature surface so that the height of the fluid passage is reduced in a downstream direction thereof.\nIn preferred embodiments, the fluid passage is structured such that the cross sections of the nozzle blades at the tip and the root of the nozzle blades are shifted in the circumferential direction of the annular fluid passage.\nA throat width between adjacent two nozzle blades is determined by the relationships:\nSpxe2x89xa6Sr less than St\nwhere Sp is the throat width at the central portion in the lengthwise direction of the nozzle blades, Sr is that at the root and St is that at the tip.\nIn other side of the present invention, the axial-flow turbine comprising: an outer diaphragm ring and an inner diaphragm ring forming together an annular fluid passage; and a plurality of nozzle blades disposed in the annular passage,\nwherein the height of the nozzle blades at a trailing edge is made to be larger than the height of the nozzle blades at a leading edge,\nthe annular fluid passage having a stepped portion at an inner surface of the outer diaphragm ring, the stepped portion having a curvature surface so that the height of the fluid passage is reduced in a downstream direction thereof, the height of the fluid passage being enlarged at a position adjacent to he trailing edge of the nozzle blades,\nthe inner trailing edge of the nozzle blades being positioned in the most downstream position and the outer trailing edge being positioned in the most upstream position.\nIn preferred embodiments, the stepped portion has a height in a radial direction of the fluid passage, the height being described by the relationships:\n0.1 less than h2/L1 less than 0.2\nwhere L1 is the height of a leading edge of the nozzle blades and h2 is te height of the stepped portion provided for the outer diaphragm ring.\nThe above object can be achieved according to the present invention by providing an axial-flow turbine comprising: a rotor wheel; a plurality of moving blades disposed on an outer surface of the rotor wheel; and an annular cover attached to a tip each of the moving blades, the annular cover and the rotor wheel forming a annular fluid passage,\nwherein the moving blades are formed into a warped shape in such a manner that the lengthwise directional central portion of the moving blades at the trailing edge of the moving blades is lower than a straight line connecting an trailing edge at the root and a trailing edge at the tip to each other.\nIn preferred embodiments, the annular fluid passage has a stepped portion at an outer surface of the rotor wheel and an inner surface of the cover, the stepped portion having a curvature surface so that the height of the fluid passage is reduced in a downstream direction thereof.\nThe stepped portion has a height in a radial direction of the fluid passage, the height being described by the relationships:\n0xe2x89xa6h3/L3 less than 0.05\n0.1 less than h4/L3 less than 0.2\nwhere L3 is the height of the leading edge of the moving blades, L4 is the height of the trailing edge of the moving blades, h3 is the height of the stepped portion provided for the rotor wheel and h4 is the height of the stepped portion provided for the cover.\nEach of said moving blades is structured such that an axial distance from points on a line connecting an trailing edge of a root of the moving blades to a trailing edge at the tip to points on a curved line forming a trailing edge of the moving blades are longest in the central portion of the lengthwise direction of the moving blades at the trailing edge of the moving blades.\nThe fluid passage is structured such that the inner surface of the cover and the outer surface of the rotor wheel are outwards inclined in the downstream direction.\nAn angle of inclination of said fluid passage is descried by the relationships:\n0xc2x0xe2x89xa6xcex81 less than xcex83 less than xcex82\nwhere xcex81 is an angle of inclination of the outer surface of the rotor wheel, xcex82 is an angle of inclination of the inner surface of the cover at the leading edge of the moving blades and xcex83 is an angle of inclination of a portion of the inner surface of the cover following the trailing edge of the moving blades.\nThe height L4 of the moving blades at the trailing edge is made to be larger than the height L3 of the moving blades at the leading edge (L3xe2x89xa6L4).\nThe fluid passage is structured such that the inner surface of the cover is outwards inclined in the downstream direction and the outer surface of the rotor wheel is inwards inclined in the downstream direction.\nAn angle of inclination of said fluid passage is descried by the relationships:\nxcex81 less than 0xc2x0 less than xcex83 less than xcex82\nwhere xcex81 is an angle of inclination of the outer surface of the rotor wheel is, xcex82 is an angle of inclination of the inner surface of the cover at the leading edge of the moving blades and xcex83 is an angle of inclination of a portion of the inner surface of the cover following the trailing edge of the moving blades.\nThe fluid passage is structured such that the cross sections of the outer and inner portions of the moving blades are shifted in the circumferential direction of the rotor wheel.\nA throat width between adjacent two moving blades is determined by the relationships:\nSr greater than Sp less than St\nwhere Sp is the width of the central portion in the lengthwise direction of the moving blades, Sr is that at the root and St is that at the tip.\nIn other side of the present invention, an axial-flow turbine comprising: a rotor wheel; a plurality of moving blades disposed on an outer surface of the rotor wheel; and an annular cover attached to an outer end each of the moving blades, the annular cover and the rotor wheel forming a annular fluid passage,\nwherein said annular fluid passage has a stepped portion at an outer surface of the rotor wheel and an inner surface of the cover, the stepped portion having a curvature surface so that the height of the fluid passage is reduced in a downstream direction thereof.\nIn preferred embodiments, the fluid passage is structured such that the cross sections of the outer and inner portions of the moving blades are shifted in the circumferential direction of the rotor wheel.\nA throat width between adjacent two moving blades is determined by the relationships:\nSr greater than Sp less than St\nwhere Sp is the width of the central portion in the lengthwise direction of the moving blades, Sr is that at the root and St is that at the tip.\nThe turbine nozzle or the turbine moving blades having the above-mentioned structure of the present invention causes the working fluid introduced to the portions adjacent to the tip and root portions of the fluid passage by the nozzle blade or the moving blade to be narrowed by the stage of the wall of the fluid passage. Thus, eddies in the secondary flow between blades can be prevented and the secondary loss can be reduced.\nSince the trailing edge of the nozzle blade and that of the moving blade are disposed downstream in the central portion of the blade, the flow lines of the portions of the working fluid in each of the nozzle blade and the moving blade are shifted to the tip and root portions. As a result, the distribution of the flows in the lengthwise direction of the blade can be uniformed. Thus, energy can effectively be converted by the moving blades. Therefore, the above-mentioned functions improve the performances of the turbine stages."}
{"text": "The placement of synthetic meshes into the body of a patient has become routine in the field of surgical practice. The most common use for mesh implants is the placement into the abdominal cavity for the repair of ventral and incisional hernias. The meshes are either placed extraperitoneally in onlay or sublay technique or intraperitoneally in intraperitoneal onlay mesh technique (IPOM). Inlay technique also played a role in the past. IPOM technique is indicated in multiple defect hernias (also known as Swiss cheese hernias) as well as for the treatment of large abdominal wall defects with loss of domain in obese patients. The implants are usually fixed by sutures, tacks or anchors. Complications found with IPOM technique are seroma formation, impaired functionality of the abdominal wall, chronic pain and adhesion formation. Chronic pain (continuous or intermittent pain persisting for more than 3 months postoperatively is mainly caused by perforating fixation devices with a reported incidence of 5-15% in open and laparoscopic hernia repair operations.\nClinical and experimental studies have shown that synthetic meshes as well as fixation devices can evoke a foreign body reaction. Unspecific tissue inflammation provokes tissue adhesions often elicited by protruding tacks, sutures or sharp mesh margins. These adhesions are potentially resulting in major complications like bowel adhesion, bowel obstruction or perforation.\nDifferent mesh concepts for adhesion prevention have been developed including coated meshes, developed for separation of peritoneal defects and used for the individual coating of meshes. Large pores and high flexibility increase mesh integration into the abdominal wall and provide good biomechanical function. Examples for coated meshes are Parietex Composite®, Sepramesh® and Proceed®. The idea of integrating mesh and antiadhesive layer in the implant is to separate implant and viscera until the mesh is covered by neomesothelium (after approximately 10 days) and in the following to reduce the foreign body reaction and adhesion formation triggered by the implant.\nHowever, using coated meshes may still trigger adverse reaction due to the use of mechanical fixation means such as sutures, tacks and anchors, which may extend from the implant. It is therefore a goal of the present invention to provide improved means to prevent or minimize all causes of unwanted tissue adhesion to the implant or its fixtures."}
{"text": "The so-called MIMO technology (multiple-input and multiple-output), which may be used in wireless radio frequency telecommunications, relates to the use of multiple send and receive antennas for a wireless communication at for example a base station or a user equipment. The MIMO technology forms the basis for coding methods which use the temporal as well as the spatial dimension for transmitting information enabling therefore a space and time coding. Thus, a quality and data rate of the wireless communication may be increased.\nIn a so-called massive MIMO system a plurality of user equipments may be arranged within a cell served by a base station having a plurality of antennas. In the massive MIMO system a configuration of individual antenna transceivers of the base station may vary depending on a location of each of the user equipments and transmission conditions in an environment of the base station and the user equipment.\nThe massive MIMO system may be used in connection with a time division duplex (TDD) system in which a transmission of an information stream between the base station and a user equipment is split up in time slots. Different time slots for uplink (UL) communication and downlink (DL) communication may be provided for communicating information from the user equipment up to the base station and for communicating information from the base station down to the user equipment. For such a combined massive MIMO TDD system there is a need for an additional time slot which may be called “header”. The header is used to send a training signal or training sequence from the user equipment to the base station. Based on the received training sequence, the base station may configure the transceivers of its antenna array. Thus, high antenna gain for the payload transmitted in the following time slots can be achieved. The payload may be transmitted in a number of uplink and downlink time slots. However, when the user equipment is moving, the channel quality may degrade due to a change of the spatial arrangement between the base station and the user equipment. The faster the user equipment is moving, the faster the channel quality may degrade. The channel quality may be measured as a bit error rate (BER). If there are multiple UL/DL slots within one frame, depending on the moving speed of the user equipment, only the first time slots of the frame may be adapted to be used before being degraded significantly by the increased bit error rate. To avoid this problem, the header rate may be increased such that the training sequence is transmitted more frequent. However, this will have impact on the system capacity, especially as all user equipments in the cell and in neighboring cells need to be synchronized.\nTherefore, there is a need to improve a data transmission in a massive MIMO system to mitigate deterioration due to mobility of a user equipment."}
{"text": "1. Field of the Invention\nThe invention relates to a device for transporting people by a spreader having a vertically movable head beam with horizontally extendable long bars.\n2. Description of the Related Art\nWhen unloading containers transported on ships, at first the containers have to be unshipped which are positioned on the hatch covers. In order to protect them against falling off during the journey the containers are locked against each other and relative to the hatch covers by means of so-called “twistlocks” (interlocking elements to be provided at the corners of the containers). Prior to the unloading procedure, said twistlocks have to be unlocked manually. In the case of twistlocks which must be handled completely manually, they are put first, also for the loading procedure, on the hatch cover or on the containers already loaded, whereupon they are also manually interlocked. As the containers are often stacked up to five layers on the hatch covers, it is extremely difficult to reach them from the deck or from the gangways between the container stacks. The lower two layers are locked or unlocked by rods from the deck or the gangway, whereas the upper layers have to be unlocked from above, using the afore mentioned rods. As in many cases salient containers hinder accessibility, the before mentioned operations are extremely dangerous and often result in fatal accidents. This is why in many States legal provisions have been issued generally prohibiting any operation from the upper surface of the containers. However, said provisions cannot be followed in practice because, up to now, suitable devices are missing which could safely transport the staff on the spot to open the twistlocks.\nIn view of these problems various producers are trying to find a solution and offer appropriate devices. Said devices (so-called personnel cages) are used to transport the staff as near as possible to their point of intervention. Personal cages are known which—when hanging under the spreader—are positioned on the uppermost container. From here, the staff operates using the already mentioned rods. In a further known device, two cages occupied by operators and also hanging underneath the spreader, hang down along the front parts of the containers, such that they make it possible to safely reach at least the twistlocks of the uppermost layer. A further known device facilitates, by means of three personnel cages interconnected through ladders, the locking and unlocking of the three uppermost layers. However, this device, due to its overall height, is only fit for limited use (necessary hoisting height and salient containers). A further important requirement to facilitate an efficient routine of loading and unloading is the necessity to mount said device directly on the bridge crane use for ship loading. Due to the specific construction of these very expensive bulk devices there is no place suitable for such a device."}
{"text": "Operation of turbine engines is well known and involves the use of a combustor in order to combust a fuel and therefore provide rotary power for the prime mover engine. Turbine engines are used in a wide variety of situations and the present invention has particular applicability to stationary engines utilised for electrical power generation and providing power for pumping operations etc. It will be understood that turbine engines are utilised in situations such as oil platforms where a local electricity generator and possibly pumping operations are required. In these situations generally natural gas drawn from production will be used in normal operation. However, during construction of the platform and during periods of no gas production the turbine engine must still remain operational at least to provide an electrical power source. Of particular advantage in situations such as oil platforms where there is little space is the high thermal efficiency of turbine engine with low emissions and package size. In such circumstances turbine engines are generally optimized for gas combustion and it is difficult to provide dual fuel operations due to the limited package space available for the fuel injection arrangement as well as costs if intrusive techniques are required to convert between fuels."}
{"text": "A concrete slab is a common structural element in buildings, where slabs of steel-reinforced concrete, e.g., between about 4 and about 20 inches thick, are used to construct floors and ceilings in buildings. For example, in many domestic and industrial buildings, a thick concrete slab (e.g., supported on foundations or directly on the subsoil) is used to construct the ground floor of a building. Thinner slabs of concrete (e.g., between about 2 and about 6 inches) may be used for exterior pavings and the like. Also, in high-rise buildings and skyscrapers, thinner concrete slabs may form the floors and ceilings on each level.\nA concrete slab may be prefabricated or poured in situ. In situ concrete slabs are typically built on a construction site using formwork (also referred to herein as a “concrete form”), which may include a type of boxing or framing into which wet concrete is poured. If a concrete slab is to be reinforced, reinforcing bars (commonly referred to as “rebars”) may be positioned within the formwork before the concrete is poured. Plastic tipped metal, or plastic bar chairs may be used to hold the rebar away from the bottom and sides of the formwork, so that, when the concrete sets, the concrete completely envelops the reinforcements. For a ground slab, the formwork may include sidewalls pushed into the ground or otherwise braced/supported. For a suspended slab, the formwork may be shaped like a tray, often supported by a temporary scaffold or the like until the concrete sets. The concrete form may be built using wooden planks and boards, plastic, steel, and the like.\nThe concrete form may use “kickers” or the like to hold the form in place, i.e., to stabilize or brace the concrete form. Kickers may include a wood block or board that is attached to the outside of a formwork member in a concrete form to provide stability to the structure. In formwork, a kicker may act as a haunch to take the thrust of another member. Kickers are often used once, or a relatively limited number of times, and then discarded. And because kickers often include wooden blocks or boards that could have otherwise been used in construction, discarding kickers after one use may waste valuable building material.\nThere remains a need for improved devices, systems, and methods for stabilizing a concrete form or other structure."}
{"text": "Pseudaletia unipunctra, the true armyworm, is endemic on all continents, being an indiscriminate pest of field crops and horticultural species. The moth is known to be migratory and in Canada (and some other temperate regions) the species enters afresh each spring from the south. Because of this behaviour, the pest may appear suddenly in large numbers. The moths lay eggs in May/June which hatch into the damaging larval (armyworm) stage during June/July. There are two generations per year in most places. Each generation consists of an egg stage, which hatches into ti larvae, which when fully grown, enter the pupal stage, during which metamorphosis occurs giving rise to the adult moth. These then mate and lay eggs so that the cycle begins over again. Thus, if moths are present in a locale, eggs and armyworms can be expected to appear there subsequently.\nThe presence of the moth in a locality is usually assessed by light trapping, wherein a bright light is used to lure the moths to a trap at night, during the periods when the moth is in flight. This technique has long been known and remains in wide use today. However, light lures indiscriminately bring many other moths, along with beetles and other insects, to the trap also; sorting is thus necessary. In practice, comparative light trapping is expensive and time-consuming, and expertise is required. A more convenient method of sampling local populations would use a simpler, cheap trap and a specific lure.\nThe compound (Z)-11-hexadecen-1-yl acetate (Z11-16:Ac) has been found in abdominal tip extracts and in effluvium from female moths of the common armyworm and has been confirmed as the primary sex pheromone component for P. unipuncta (see Hill and Roelofs \"Environmental Entomology\", Vol. 9, No. 4, p. 408-411, August 1980). McDonough et al, J. of Chem. Ecology, Vol. 6, No. 3, p. 565-572, 1980, also identified Z11-16:Ac, a hexadecen-1-o1 (believed to be (Z)-11 isomer), and possible (Z)-9-hexadecen-1-yl acetate, in extracts of the female armyworm moths. In field tests, no combination of compounds was found to be more attractive than Z11-16:Ac alone. Farine et al confirmed the presence of Z11-16:OH in the natural extract and found (in addition to Z11-16:Ac) that hexadecanyl acetate and Z9-16:Ac were also present (see C.R. Acad. Sc. Paris, t. 292, Series III, p. 101-104, Jan. 5, 1981).\nIn early tests, we had found that Z11-16:Ac with 0.1-0.5% (Z)-11-hexadecenol produced a reasonably effective and specific lure for the male armyworm moths. The attraction was increased on adding trace amounts (e.g. 0.01%) of (Z)-9-tetradecen-1-yl acetate (Z9-14:Ac). (See W. Steck et al, \"Environmental Entomology\", Vol. 9, No. 5, p. 583-585, October 1980). In tests on Polia atlantica we had tried the ternary mixture Z11-16:Ac+Z11-16:OH+Z11-16:Ald in the weight ratios 200:1:1 and found it ineffective for this species (see Table 2 in this latter reference). Some or all of these latter three compounds are known to be present in attractants for other species but the ratios are vastly different and the previously described mixtures are substantially nonattractive to armyworm moths."}
{"text": "1. Technical Field\nThe present invention relates, generally, to a control system for maintaining optimum efficiency of a backlight and, more particularly in a preferred embodiment, to a closed loop temperature controller for adjusting the temperature within a fluorescent lamp to thereby optimize lamp arc drive for a given predetermined brightness set point.\n2. Background Art and Technical Problems\nScreen displays which employ fluorescent lamp backlights are used extensively in commercial, military, and consumer electronic applications. For example, such backlights are commonly used in desktop computers, laptop computers, screen displays for industrial equipment, and in connection with \"heads up\" or other screen displays in the cockpits in both commercial and military aircraft.\nConventional fluorescent lamps are commonly employed in backlit Liquid Crystal Display (LCD) applications. In a typical LCD, alphanumeric characters and other graphical images are produced on the viewing screen by selectively energizing or de-energizing preselected pixels in a two dimensional matrix to display the information. In a normally black screen display, predetermined pixels are illuminated to display the data or information as illuminated characters on a black (or other dark shade) background. In a normally white display, on the other hand, the desired data and/or information corresponds to the non-illuminated pixels, such that the information appears as black (or other dark color) images on a white (or other light color) background. In either case, a bright, consistent \"background\" light is necessary to achieve desirable contrast on the flat screen display. Indeed, in certain applications (e.g., military avionics), the high contrast provided by a bright backlight is essential to proper operation of the display.\nIt is also desirable to obtain a desired brightness while minimizing power consumption. This is particularly important in portable electronics, for example laptop computers and the like, where battery life is an important product feature.\nPresently known systems for controlling the brightness of a fluorescent backlight lamp typically involve a control system for supplying lamp arc drive to the backlight, to thereby excite the gas atoms within the sealed lamp enclosure to create visible light. The amount of visible light emitted by the lamp is sensed, for example by a photodiode, and a feedback signal indicative of the brightness output of the lamp is fed back to a control circuit. This feedback signal (indicative of actual brightness) is compared to an input signal representative of a desired brightness level, and presently known control systems drive the difference between this actual signal and the desired signal to a minimum. Under this control regime, if the actual brightness is less than the desired brightness, the controller increases the lamp arc drive applied to the lamp until the actual brightness equals the desired brightness. If, on the other hand, the actual brightness is greater than the desired brightness, the controller circuit reduces the magnitude of the lamp arc drive applied to the lamp until the actual brightness emitted from the lamp again equals the desired brightness level for the lamp. Presently known prior art brightness control systems typically employ a \"cold spot\" at a predetermined point on the lamp which functions to keep a certain amount of the gas (typically mercury) within the lamp in a condensed state. Such \"cold spot\" systems employ the well known principle that maintaining the temperature of the cold spot in a specified range allows for very efficient operation of the lamp. Presently known systems, however, often require expensive components to maintain the cold spot, and do not adequately compensate for drifting or degradation over time of some of the parameters which influence the efficiency of the lamp.\nA fluorescent lamp control system is thus needed which overcomes the shortcomings of the prior art."}
{"text": "The present invention relates to a tip turbine engine, and more particularly to mounting the tip turbine engine along a single mounting plane.\nAn aircraft gas turbine engine of the conventional turbofan type generally includes a forward bypass fan and a low pressure compressor, a middle core engine including a combustor, and an aft low pressure turbine all located along a common longitudinal axis. Although highly efficient, conventional turbofan engines operate in an axial flow relationship. The axial flow relationship results in a relatively complicated elongated engine structure and therefore requires multiple mounting planes along the elongated engine structure to support the engine and mount the engine on an aircraft. Utilizing multiple mounting planes may complicate the mounting process, add weight, and make assembly laborious and expensive.\nA recent development in gas turbine engines is the more longitudinally compact tip turbine engine. Tip turbine engines locate an axial compressor forward of a bypass fan, which includes hollow fan blades that receive airflow from the axial compressor therethrough such that the hollow fan blades operate as a centrifugal compressor. Compressed core airflow from the hollow fan blades is mixed with fuel in an annular combustor located radially outward from the fan. The combustor ignites the fuel mixture to form a high energy gas stream which drives turbine blades that are integrated onto the tips of the hollow bypass fan blades for rotation therewith as disclosed in U.S. Patent Application Publication Nos.: 2003192303; 20030192304; and 20040025490. The tip turbine engine provides a thrust to weight ratio equivalent to conventional turbofan engines of the same class within a package of significantly shorter longitudinal length.\nAccordingly and because of the unique architecture and shorter longitudinal length of the tip turbine engine, it is desirable to provide more efficient and simplified mounting assemblies for the tip turbine engine."}
{"text": "The present invention relates to a document with anti-counterfeiting means to prevent reproducibility by photocopying.\nIt is known that photocopying machines capable of reproducing any document in the finest of details are currently commercially available. These machines have already been used for the fraudulent duplication of documents, such as for example checks and the like, obtaining a photocopy that is substantially identical to the original and cannot therefore be distinguished easily.\nIn order to solve this problem, various methods have already been marketed which, for example, produce on the document regions that are not visible within the normal spectrum of visible light but instead appear in the photocopy, thus making the counterfeiting evident.\nHowever, these methods entail considerable complexity during production of the document and accordingly so far they have not been widely adopted."}
{"text": "The present invention relates to charge forming systems, and more particularly to a system which automatically and continuously seeks an optimum air/fuel mixture.\nThe provision of an instantaneously correct fuel/air mixture to an internal combustion engine over a full range of speeds and loads has for years occupied countless engineers, technicians and inventors. Accordingly, the charge forming apparatus have progressed over the years from crude mixing devices to complex, sophisticated carburetion and fuel injection systems. Nonetheless, it is recognized that to date no charge forming system has been implemented which will provide the ideal fuel/air mixture to an internal combustion engine over a complete operating range.\nOf the various approaches to charge forming systems, two basic lines of attack may be discerned. With one approach, air inflow is controlled by a throttle or the like and fuel is caused to be entrained with the passing air, in response to the air flow characteristics. The metering of the mixing apparatus is such that a relatively constant, predictable mixture results over the anticipated air flow range. This is the basic rule of operation of the carburetor, which is of course the predominant charge forming mechanism used with internal combustion engines. The same thesis, however, has also been implemented in many forms of fuel injection apparatus.\nThe second approach referred to above may be referred to as a \"programmed\" approach inasmuch as the charge forming system is \"programmed\" or adjusted to respond to certain stimuli to effect a desired fuel/air ratio. The stimuli applied to such systems may include temperature, pressure, engine speed, and the like. Conventionally, such systems make use of not one but a plurality of stimuli in order to more accurately react to inferred engine operating conditions. Charge forming systems of this type thus react to given types of stimuli in a predetermined manner, effecting a fuel/air ratio which is assumed to be the \"correct\" one for the engine. In theory this approach should work well, assuming that a great number of preprogrammed responses are available to correspond to the almost limitless combinations of stimuli which occur over the full range of operation for most internal combustion engines, especially those used in vehicles.\nIn recent years the increased recognition of environmental degradation caused by internal combustion engine exhaust, along with the need for economical operation, has caused renewed interest in more efficient and sophisticated charge forming control systems. Generally, the approach taken in developing more responsive charge control systems has been to make use of the most modern and sophisticated technology, particularly in the field of electronics, to accommodate the functional complexities which are inherent in the \"programmed\" type of system. Accordingly, control systems which are in effect small analog or digital computers are created which can \"tailor\" a fuel/air ratio to a given operating condition, as evidenced by various stimuli. The stimuli or sensed parameters include exhaust gas temperature and composition, engine temperature, combustion pressure, inlet manifold and/or venturi vacuum, engine speed, along with many other operating parameters. However, inasmuch as these systems all provide a fixed, predetermined response to a given set of stimuli, they cannot be flexible enough to find the true optimum fuel/air ratio for all operating conditions, but rather provide mixture ratios which are only assumed to be correct.\nAccordingly, it will be appreciated that it would be desirable to provide an improved charge forming system which consistently achieves the actual optimum fuel/air ratio for any given operating condition.\nIt is therefore an object of the present invention to provide an improved charge forming system for optimizing the mixture ingested by an internal combustion engine over a broad operating range.\nIt is another object of the present invention to provide an improved charge forming system which does not provide a predetermined or programmed mixture to an internal combustion engine.\nYet another object is to provide an improved fuel metering system for supplying vaporized fuel to an internal combustion engine.\nIt is a further object to provide a charge forming system for an internal combustion engine which effects a marked reduction in the output of pollutants by the engine."}
{"text": "[Metallocene Compounds]\nIn recent years, metallocene compounds are well known as homogeneous catalysts for olefin polymerization. After the report of isotactic polymerization by W. Kaminsky et al. (see Non Patent Literature 1), many studies have been made on olefin polymerization, in particular, stereoregular α-olefin polymerization using metallocene compounds.\nIn α-olefin polymerization using metallocene compounds, it is known that the stereoregularity and the molecular weights of the obtainable α-olefin polymers are greatly varied by the introduction of substituents into the cyclopentadienyl ring ligands of the metallocene compounds or by the bridging of the two cyclopentadienyl rings.\n[Bridged Metallocene Compounds]\nFor example, there are the following reports as to propylene polymerization catalyzed by metallocene compounds which have a ligand in which a cyclopentadienyl ring and a fluorenyl ring is bridged to each other.\nFrom the viewpoint of stereoregularity, dimethylmethylene(cyclopentadienyl) (fluorenyl)zirconium dichloride affords syndiotactic polypropylene (see Non Patent Literature 2); dimethylmethylene(3-methylcyclopentadienyl) (fluorenyl)zirconium dichloride having a methyl group at the 3-position of the cyclopentadienyl ring affords hemiisotactic polypropylene (see Patent Literature 1); and dimethylmethylene(3-tert-butylcyclopentadienyl)(fluorenyl) zirconium dichloride having a tert-butyl group at the 3-position of the cyclopentadienyl ring affords isotactic polypropylene (see Patent Literature 2).\nDimethylmethylene(3-tert-butyl-5-methylcyclopentadienyl)(3,6-di-tert-butylfluorenyl)zirconium dichloride having tert-butyl groups at the 3- and 6-positions of the fluorenyl ring affords polypropylene with higher isotactic stereoregularity than obtained with dimethylmethylene(3-tert-butyl-5-methylcyclopentadienyl) (fluorenyl)zirconium dichloride (see Patent Literature 3).\nFrom the viewpoint of molecular weights, diphenylmethylene(cyclopentadienyl) (fluorenyl)zirconium dichloride having a cyclopentadienyl ring and a fluorenyl ring bridged via diphenylmethylene affords syndiotactic polypropylene having a higher molecular weight than obtained with dimethylmethylene(cyclopentadienyl) (fluorenyl)zirconium dichloride (see Patent Literature 4); diphenylmethylene(3-(2-adamantyl)-cyclopentadienyl)(fluorenyl)zirconium dichloride having a diphenylmethylene bridge affords isotactic-hemiisotactic polypropylene having a higher molecular weight than obtained with dimethylmethylene(3-(2-adamantyl)-cyclopentadienyl)(fluorenyl)zirconium dichloride (see Non Patent Literature 3); and dimethylmethylene(3-tert-butyl-5-methylcyclopentadienyl) (fluorenyl)zirconium dichloride having a methyl group at the 5-position of the cyclopentadienyl ring (the α-position relative to the bridge) affords isotactic polypropylene having a higher molecular weight than obtained with dimethylmethylene(3-tert-butylcyclopentadienyl)(fluorenyl) zirconium dichloride (see Patent Literature 5).\nFurther, dimethylmethylene(3-tert-butyl-2-methylcyclopentadienyl) (fluorenyl)zirconium dichloride and diphenylmethylene(3,4-dimethylcyclopentadienyl) (fluorenyl) zirconium dichloride having substituents at two adjacent positions on the cyclopentadienyl ring afford polypropylene having a lower molecular weight than obtained with dimethylmethylene(3-tert-butyl-5-methylcyclopentadienyl) (fluorenyl)zirconium dichloride and diphenylmethylene(3-methylcyclopentadienyl) (fluorenyl)zirconium dichloride, respectively (see Patent Literatures 5 and 6).\n[5-Membered Ring-Bridged Metallocene Compounds]\nA study reports the polymerization of propylene catalyzed by a metallocene compound in which a cyclopentadienyl ring and a fluorenyl ring are bridged via a 5-membered ring. However, such metallocene compounds have low usefulness in industry because of the fact that the stereoregularity of the obtainable polypropylenes is very low (see Non Patent Literature 4).\nA recent study reports a metallocene compound having a cyclopentadienyl ring and a fluorenyl ring bridged via a 5-membered ring which can afford polypropylene having relatively high stereoregularity (see Patent Literature 7).\nThese metallocene compounds mentioned above exhibit excellent polymerization performance. In some applications, however, the catalysts are often required to afford polymers having still higher stereoregularity or still higher molecular weight with higher economic efficiency, namely, with high catalytic activity even under high-temperature polymerization conditions. Improvements are thus required.\n[Metallocene Compounds Having Substituted Indenyl Ligands]\nAccording to reports, metallocene compounds having substituted indenyl ligands afford relatively high stereoregularity or molecular weight (see Patent Literatures 8 and 9). However, such compounds are unsatisfactory in terms of performance under economically efficient polymerization conditions.\nBecause metallocene compounds are soluble in reaction media, they are generally used to catalyze polymerization in the form of supported catalyst systems in slurry polymerization or gas phase polymerization. Specifically, the metallocene compounds are supported on solid carriers. However, it is known that the polymerization performances such as stereoregularity control of the aforementioned compounds are markedly decreased when they are used in the supported form on carriers as compared to in the absence of carriers.\n[Metallocene Compounds Having Substituted Azulenyl Groups]\nTo solve such problems, for example, a recent study reports a metallocene compound having a substituted azulenyl group as a ligand (see Patent Literature 10). However, even such catalysts do not achieve sufficient performances such as stereoregularity control when the polymerization temperature is elevated to obtain economic efficiency or when the compounds are supported on solid carriers.\nUnder these circumstances, there have been demands for further improvements in the catalytic performances such as polymerization activity, stereoregularity control and molecular weight control of polymerization catalysts including metallocene compounds (hereinafter, also written as “metallocene catalysts”).\n[Reports of Polymerization Using Metallocene Catalysts]\nRegarding the polymerization of monomers other than propylene with use of metallocene catalysts, for example, there are studies reporting the polymerization of 1-butene catalyzed by a metallocene compound having an indene ring as a ligand (see Patent Literatures 11 and 12).\nFurther, the polymerization of 1-butene catalyzed by a metallocene compound having a fluorene ring as a ligand has been reported. Specifically, the polymerization using isopropylidene (3-t-butyl-5-methylcyclopentadienyl) (fluorenyl)zirconium dichloride has been disclosed (see Patent Literature 13). According to this report, the obtainable polybutene has lower regioirregularity due to 4,1-insertions and exhibits higher heat resistance and mechanical strength as compared to when the polymerization is catalyzed by a metallocene compound having an indene ring as a ligand. In some applications, however, the catalysts are often required to afford polymers having still higher stereoregularity or still higher molecular weight under more economically efficient conditions. Improvements are thus required.\nA study reports the polymerization of 4-methyl-1-pentene as a main monomer (see Patent Literature 14). Patent Literature 14 describes that metallocene-catalyzed polymers outperform conventional Ziegler-Natta-catalyzed polymers in the balance of properties such as heat resistance and are thus highly useful in industry. However, applications sometimes require that the polymers have a still higher melting point/stereoregularity or have a still higher molecular weight.\nIn general, increasing the polymerization temperature enhances economic efficiency but at the same time tends to decrease the molecular weight or the melting point/stereoregularity of the obtainable polymers. Thus, there have been demands for the development of novel polymerization processes capable of producing polymers having a higher molecular weight and a higher melting point/higher stereoregularity."}
{"text": "In ophthalmic surgery such as treatment of cataract, corneal transplantation, or vitreous surgery, an intraocular irrigating solution is used to protect the intraocular tissue during the surgery.\nAn intraocular irrigating solution containing a glutathione derivative has been developed as such an intraocular irrigating solution and has been considered to be a very stable intraocular irrigating solution that effectively suppresses precipitation of calcium (see Patent Literature 1).\nMoreover, an intraocular irrigating solution containing a phosphoric diester of ascorbic acid and the tocopherol has also been developed and considered to be excellent in protection of corneal cells and to enable ophthalmic surgery to be safely performed (see Patent Literature 2).\nHowever, cysteine persulfide (CysSSH) as an active sulfur molecule has been reported to be produced in vivo. For example, pyridoxal phosphate (PLP) well known as a coenzyme such as a deaminase has been reported to directly react with cystine to produce CysSSH. Moreover, a cysteine residue in the active center of a certain enzyme, particularly rhodanese known as an enzyme for detoxifying cyan (CN−), has been known to receive a sulfur molecule supplied from a thiosulfate ion to form a persulfide, which reacts with a cyanide ion and subjected to metabolic detoxification into a thiocyanate ion (SCN−).\nCommonly, CysSSH had been recognized as an unstable metabolic intermediate responsible for the catalytic activity of such a special enzyme. However, the present inventor found that an active sulfur molecule such as CysSSH in vivo exhibits high antioxidant activity and the function of controlling a redox signal (see Non Patent Literature 1)."}
{"text": "The present application relates to a non-planar semiconductor structure and a method of forming the same. More particularly, the present application relates to a semiconductor structure containing tensile strained silicon-containing fins for n-type FETs and compressive strained silicon germanium alloy fins for p-type FETs.\nFor more than three decades, the continued miniaturization of metal oxide semiconductor field effect transistors (MOSFETs) has driven the worldwide semiconductor industry. Various showstoppers to continued scaling have been predicated for decades, but a history of innovation has sustained Moore's Law in spite of many challenges. However, there are growing signs today that metal oxide semiconductor transistors are beginning to reach their traditional scaling limits. Since it has become increasingly difficult to improve MOSFETs and therefore complementary metal oxide semiconductor (CMOS) performance through continued scaling, further methods for improving performance in addition to scaling have become critical.\nThe use of non-planar semiconductor devices such as, for example, semiconductor fin field effect transistors (finFETs) is the next step in the evolution of complementary metal oxide semiconductor (CMOS) devices. Semiconductor fin field effect transistors (FETs) can achieve higher drive currents with increasingly smaller dimensions as compared to conventional planar FETs. In order to extend these devices for multiple technology nodes, there is a need to boost the performance with high mobility channels and stressor elements."}
{"text": "1. Field of Invention\nA casino-type game of chance utilizing dice a game board, betting pieces and dice.\n2. Prior Art\nThere are various existing casino-type dice games. Such games often require players to select various outcomes of dice rolls by placing bets in designated areas. Although simple in nature, many of the dice games, such as craps, have complex rules and tend to attract only experienced gamblers. Therefore, there has been a need for casino-type dice games with simple rules that cater to less experienced gamblers.\nThis invention relates to a casino-type dice game. Preferably, the game utilizes two six-sided dice with color-coded indicia arranged on the faces of each dice cube so that each face has a different color. Other indicia may be used that permit each face of the cube to be distinguished from one another. Thus, the theme of the game is open to utilizing different indicia such as numbers, animals, sporting or business themes.\nThe present game is a game of chance. This is accomplished by dividing the game into three betting areas associated with different indicia combinations. Although the odds to select the correct combination are the same for all three betting areas, the stakes are modified to represent a different payout for each betting area. In addition, each betting area represents a different duration for the period a wager may remain without being lost.\nThe present invention involves a plurality of players and is played with at least two dice and a game board. Three betting areas are provided and each is associated with different indicia combinations that can exist when the dice are rolled. In addition, rewards are set for each betting area based on indicia combinations located therein.\nThe game is played by the players placing wagers next to different indicia combinations located in at least one of the three betting areas. The dice is rolled over a sequence of times. The players are rewarded after each dice roll in an amount associated with a respective betting area.\nThe three betting areas are distinguished as follows. In the first betting area, a first set of wagering spots represent different indicia combinations that correspond to five dice surfaces wherein each of said indicia combinations located in the first set of wagering spots comprise different indicia. In the second betting area, there is a second set of wagering spots representing different indicia combinations that correspond to a sixth dice surface matched with one of the five dice surfaces corresponding to the first set of wagering spots. Finally in the third betting area, there is a third set of wagering spots representing indicia combinations that have the same indicia.\nAdditional features of this invention will become apparent from the ensuing descriptions of this invention."}
{"text": "In order to support various safety services, a vehicle needs to observe surroundings thereof. To this end, the vehicle includes a surrounding observing device. For example, the vehicle includes a lidar sensor. The lidar sensor transmits a signal toward a front direction or the like and analyzes a feedback signal, thereby generating a video image. During this process, the lidar sensor divides upper and lower sides of the front thereof into predetermined layers and obtains the video image for each layer. That is, the lidar sensor according to the related art supports a multi-layer image.\nHowever, since the image generated based on the multi-layer has a significant amount of data distribution to be determined to support a specific service, it has several problems. Therefore, it is necessary to converge the same object among objects for each multi-layer. In order to solve this problem, according to the related art, the image has been processed based on a state in which the objects for each layer are overlapped. However, in the case in which the image is simply processed based on the overlapped state, it is impossible to finely distinguish between the objects. This may potentially make a risk prediction of the vehicle or avoidance difficult."}
{"text": "Liquid crystal displays (LCDs) find wide usage in applications such as air vehicle cockpits, automobile dashboards, etc. Such LCDs generally include a liquid crystal layer sandwiched between a pair of polarizers and a pair of electrodes. Voltage is selectively applied by the electrodes across the liquid crystal (LC) layer so as to selectively portray an image to the viewer by way of the front polarizer.\nFIG. 1(a) illustrates prior art liquid crystal display 3 in an air vehicle cockpit 5 occupied by pilot 7 and copilot 9. In prior art cockpit (or automobile dashboard) applications such as this, situations commonly arise where a portion of the image 13 emitted from the active matrix liquid crystal display (AMLCD) 3 is reflected at point 11 off of cockpit canopy 13 toward copilot 9, while another portion of the image is transmitted through canopy 13. Air vehicle cockpit 5 as shown in FIG. 1(a) is designed so as to house copilot 9 directly behind and above pilot 7 as illustrated. Thus, when a portion 15 of the image 13 from LCD 3 is reflected off of canopy 13 at reflection point 11, it can be seen by copilot 9 when the copilot is looking out of the cockpit. The other portion or remainder of the image is transmitted through the canopy as shown at 14. This reflection 15 from point 11 is, of course, undesirable in view of the fact that copilot 9 typically wishes to locate targets and the like exterior the aircraft and does not wish to be interfered with by such display reflections. Needless to say, it is undesirable to have the pilot's or copilot's vision obstructed during flight operations due to display image reflections off of the canopy.\nThe above-referenced display reflection problems inside of aircraft cockpits are most severe at night when reflections 15 exceed the intensity of the light or environment outside of the cockpit, but are also present during daylight hours. The majority of undesirable reflections 15 come from the interface between the air inside of the cockpit and the material (e.g. glass or plastic) of canopy 13. Typically, the index of refraction of cockpit canopy 13 is about 1.5. Depending upon the shape of the canopy, reflections 15 may be concentrated in convex curved areas or mirror-like in flat canopy areas.\nIn view of the above, it is clear that there exists a need in the art for a liquid crystal display and method of implementation for reducing display reflections off of exterior mediums such as canopies of aeronautic cockpits.\nFIG. 1(b) is a reflectance versus incident angle .THETA..sub.i graph from \"Optics\" by Eugene Hecht, 1987. As shown in FIG. 1(b), for an incoming unpolarized wave made up of two incoherent orthogonal polarization states, only the component polarized normal (R.sub.perpendicular) to the incident plane and therefore parallel to the surface of the medium will be reflected. The incident plane is defined by the line connecting the reflection point on the exterior medium and the display front surface and the medium surface normal The component (R.sub.parallel) parallel to the plane of incidence is not reflected and is transmitted through the medium as shown at 14. The particular angle of incidence for which this situation occurs is designated by .THETA..sub.p and is referred to in the art as the polarization angle or Brewster's angle, whereupon .THETA..sub.p +.THETA..sub.t =90.degree.. Therefore, from Snell's Law, EQU n.sub.i Sin.THETA..sub.p =n.sub.t Sin.THETA..sub.t\nand the fact that .THETA..sub.t =90.degree.-.THETA..sub.p, it follows that EQU n.sub.i Sin.THETA..sub.p =n.sub.t Cos.THETA..sub.p\nand EQU Tan.THETA..sub.p =n.sub.t /n.sub.i.\nThis is known as Brewster's Law. The above parameters and equations are defined in \"Optics\", by Eugene Hecht, the disclosure of which is incorporated herein by reference.\nThus, when the incident light ray or beam is in air (n.sub.i =1) and if the transmitting medium (i.e. exterior medium such as canopy 13 or the like) is glass, in which case n.sub.t =1.5, the polarization angle .THETA..sub.p is about 56.degree.. Similarly, this concept may be exemplified by the situation when an unpolarized beam strikes the surface of a pond (n.sub.t =1.33 for water) at an angle of 53.degree., the reflected beam will be completely polarized with its E-field perpendicular to the plane of incidence or, if you like, parallel to the water's surface.\nFor linearly polarized light having its E-field parallel to the plane of incidence, the \"amplitude reflection coefficient\" is defined as r.sub.parallel =(E.sub.or /E.sub.oi).sub.parallel, that is, the ratio of the reflected to incident electric field amplitudes. Similarly, when the electric field is normal to the incident plane, r.sub.perpendicular =(E.sub.or /E.sub.oi).sub.perpendicular. The corresponding irradiance ratio (the incident and reflected beams have the same cross-sectional area) is known as the \"reflectance\", and since irradiance is proportional to the square of the amplitude of the field, EQU R.sub.Parallel =r.sup.2.sub.parallel\nand EQU R.sub.perendicular =r.sup.2.sub.perpendicular\nSquaring the appropriate Fresnel equations results in ##EQU1##\nThe reflectance for linearly polarized light with E or a polarization direction parallel to the plane of incidence vanishes and the beam is completely transmitted when the angle of incidence .THETA..sub.i equals Brewster's angle .THETA..sub.p. This is shown in FIG. 1(b) where n.sub.t =1.5 and Brewster's angle .THETA..sub.p equals about 56.degree., so that the reflectance of R.sub.parallel is substantially zero when .THETA..sub.i substantially matches .THETA..sub.p or Brewster's angle. FIG. 1(b) is a plot for the particular case where n.sub.i =1 and n.sub.t =1.5. The middle curve corresponds to incident natural light. As shown, when the polarization direction of incoming linearly polarized light is substantially perpendicular to the plane of incidence at Brewster's angle of about 56.degree., substantial reflection results at reflection point 11. However, at Brewster's angle .THETA..sub.p of about 56.degree., when the direction of polarization of incoming light is substantially parallel to the plane of incidence (i.e. parallel to the Y-axis), the result is substantially no reflectance and nearly complete transmission at point 11. As shown by the middle curve in FIG. 1(b), when the polarization direction is mixed, reflection from point 11 results.\nFIG. 1(c) illustrates incident light 16 being reflected at point 11 off of canopy 13. Reflection 15 and transmission 14 define angles .THETA..sub.r and .THETA..sub.t with the Y-axis (normal to the surface of medium 13) respectively. The angle of incidence .THETA..sub.i is also shown in FIG. 1(c).\nU.S. Pat. No. 4,025,161 discloses a liquid crystal display device including a quarter wavelength retarder disposed on the front of the display so that the front polarizer is between the liquid crystal material and the retarder. The quarter wavelength retarder together with the front polarizer form a circular polarizer in the '161 patent. Unfortunately, if the liquid crystal display of the '161 patent were positioned in place of display 3 in FIG. 1(a), undesirable image reflection 15 would still occur in a substantial amount because the circularly polarized display image incident upon reflection point 11 would not be substantially polarized parallel to the plane of incidence and therefore reflection 15 toward copilot 9 would result thus inhibiting the copilot's view. In other words, circularly polarized light incident upon canopy 13 at reflection point 11 (even when .THETA..sub.i =.THETA..sub.p) in FIG. 1(a) would be reflected at 15 toward viewer 9, this, of course, being undesirable.\nU.S. Pat. No. 4,266,859 discloses a liquid crystal display of the reflective dichroic type, including upper and lower quarter wavelength retarders. Again, the LCD of the '859 patent suffers from the same drawbacks as those discussed above with respect to the '161 patent.\nIn view of the above, it is clear that there exists a need in the art for a liquid crystal display designed so as to reduce image reflection off of exterior mediums (e.g. automotive windshields or cockpit canopies) at a reflection point and a method of implementing same. There also exists a need in the art for a liquid crystal display which when arranged in an appropriate manner, causes the image emitted toward a reflection point on the exterior medium to have a polarization direction substantially parallel to the plane of incidence at the reflection point and .THETA..sub.i to substantially match .THETA..sub.p so that substantially no reflection results and the image is mostly transmitted through the medium instead of being reflected thereby.\nIt is a purpose of this invention to fulfill the above-described needs, as well as other needs which will become apparent to the skilled artisan from the following detailed description of this invention."}
{"text": "1. Field of the Invention\nThe present invention relates to the technical field of integrated circuits and, more particularly, to an engineering change order (ECO) hold time fixing method.\n2. Description of Related Art\nDue to a wide range of dynamic variations, e.g., supply voltage droops, process variations, temperature fluctuations, soft errors, and transistor aging degradation, the timing characterization is extremely difficult in modern IC designs. Therefore, in conventional design, designers conservatively reserve a timing guardband to ensure correct functionality even under the worst-case circumstance. However, this reserved guardband may severely degrade circuit performance, i.e., limit the clock frequency.\nAccordingly, several resilient circuits have been proposed to eliminate the guardband by error detection and correction. For example, a Razor flip-flop (FF) is used as an error detection circuit (See D. Ernst et al., “Razor: a low-power pipeline based on circuit-level timing speculation” MICRO, pp.7-18, 2003). FIG. 1 illustrates a conventional Razor FF used as the error detection circuit. As shown in FIG. 1, the error detection circuit has an extra storage element, i.e., the shadow latch 110, to sample the output of a combinational logic 120 by a delayed clock clk_dly. The output of a main flip-flop 130 is compared with the output of the shadow latch by a comparator 140, i.e., an XNOR gate. When the comparison result indicates a detected timing error, a timing error signal is generated, and the error correction is performed through an instruction replay.\nHowever, these resilient circuits require a significant hold time margin for short paths. Taking the circuit of FIG. 1 as an example, the resilient circuit may detect a false timing error if the result of the next computation is propagated through a short path and sampled by the delayed clock. To avoid such false error detection, the short paths have to exceed an error detection window w, i.e., the phase difference between the normal clock clk and the delayed clock clk_dly. The error detection window w causes an extra hold time margin requirement. This issue also exists in the resilient circuits proposed in other articles. Due to the extra hold time margin requirement, the short path padding or hold timing fixing in the resilient circuits becomes more challenging.\nTo resolve this padding problem, prior works typically insert buffers to lengthen the short paths.\nAmong the prior works, the delay padding is combined with a clock skew scheduling to minimize the clock period at the logic re-synthesis stage. The goal is to determine the padding path for each path rather than to decide where to insert the delay.\nBy contrast, another short path padding method determines the positions to insert the delay. This problem is solved by a linear programming proposed in N. V. Shenoy et al. “Minimum padding to satisfy short path constraints.” ICCAD, pp.156-161, 1993. However, such linear programming is time-consuming and not applicable to large-scale circuits. Hence, another prior art provides greedy heuristics. One greedy rule is to pad the gate with the largest setup slack for trying not to hurt the longest path delay. The other is to pad at the gate passed by the most hold violating paths for trying to reduce the total padding delay.\nFIG. 2 shows a flowchart of integrating timing error resilient circuits into a design. Because of considering the timing guardband, based on the logic synthesis and timing analysis of conservative clock period, the target clock period and the error detection window w are determined. Moreover, the invention proposes coarse-grained and fine-grained padding allocation methods to further achieve the derived padding values at physical implementation. Sth/Hth represents a ratio of the target clock period over the conservative clock period. The timing suspicious flip-flops, whose longest path delays exceed the target clock period, can be replaced by the resilient circuits. Before the replacement, the design is re-synthesized where the suspicious flip-flops are assigned with an extra hold time margin to cover the error detection window w. After the replacement, a placement and routing is applied. Because of the significant hold time margin, the hold violations may still exist in a placed and routed design. Finally, the short path padding is performed.\nIn addition, it is found that the cited greedy heuristics cannot pad the short paths well. FIGS. 3A-3D show typical short path padding. FIG. 3A gives an input design, where gates g1, g2, g3 incur hold violations. After iteratively padding delay on the gate either with the largest setup slack or with the most hold violating paths, the result is shown in FIG. 3B or 3C. From FIGS. 3B and 3C, it is found that the unresolved hold violation at gate g2 exists. However, all the hold violations can be cleaned by the padding way shown in FIG. 3D. Accordingly, it is clear that the padding based only on local views cannot pad all the short paths. Moreover, even an optimal padding solution is found, it may still fail at a physical implementation because the delay offered by buffers is fixed. For example, if one buffer offers either a 0.15-unit or 0.25-unit delay, the padding task still fails on gate g2.\nTherefore, it is desirable to provide an improved hold timing fixing method to mitigate and/or obviate the aforementioned problems."}
{"text": "The present invention relates to fluid conductivity sensors, and more particularly to a non-contacting fluid conductivity sensor which generates essentially no external current field and thus is not influenced by nearby objects in the fluid in which the sensor is immersed.\nIndustrial process control systems often require the measurement of the electrical conductivity of a fluid e.g. seawater, used in the system. Sensors for fluid conductivity measurement usually fall into one of two categories, i.e, contacting sensors and non-contacting sensors. Contacting sensors rely on a direct electrical contact between the measurement electronics and the material to the fluid, while non-contacting sensors typically employ driving and sensing transformers which, respectively, induce and measure a flow of current in the conductive fluid, the measurement of the induced current being a function of the conductivity of fluid. One problem which has existed with existing non-contacting conductivity sensors is that they may be influenced by objects a considerable distance from the sensor in the body of fluid into which the sensor is immersed. This influence is caused by the fact that conventional, immersible sensors typically have an external current field which extends a substantial distance from the sensor in the fluid in which the sensor is immersed. This is a significant problem in that calibration performed initially in the calibration laboratory will not be the same as calibration in the field."}
{"text": "Exemplary embodiments of the present invention relate to a system including a slave device and a master device.\nPackaging technology of semiconductor elements has been continuously developed according to a demand for miniaturization and high capacity. Recently, a variety of technologies for a stacked semiconductor package capable of satisfying mounting efficiency as well as the miniaturization and high capacity are being developed.\nThe stacked semiconductor package may be fabricated by the following methods. First, individual semiconductor chips may be stacked, and then packaged at once. Second, individual semiconductor packages may be stacked. The individual semiconductor chips of the stacked semiconductor package are electrically coupled through metallic wires or through silicon vias (TSV).\nHowever, in the conventional stacked semiconductor package using metallic wires, since the electrical signal exchange is performed through the metallic wires, the speed is low. Furthermore, since a large number of wires are used, electrical characteristic degradation may occur. Furthermore, since an additional area for forming the metallic wires is required in a substrate, the size of the package may increase. Furthermore, since a cap for wire bonding is required between the semiconductor chips, the height of the package may increase.\nRecently, a stacked semiconductor package using a TSV has been proposed. In general, the stacked semiconductor package is fabricated by the following method. First, a via hole is formed in a semiconductor chip so as to pass through the semiconductor chip, and a through electrode called a TSV is formed by filling the via hole with a conductive material. Then, an upper semiconductor chip and a lower semiconductor chip are electrically coupled through the through electrode.\nFIG. 1A is a block diagram illustrating a coupling state between a master device and slave devices. FIG. 1B is a diagram illustrating a state in which the slave devices are stacked and coupled to the master device.\nThe master device 100 refers to a device which controls the slave devices, and the slave devices DEV(1), DEV(2), . . . , DEV(i), DEV(j), and DEV(k) refer to devices which are controlled by the master device 100. An example of the master device 100 and the slave devices DEV(1), DEV(2), . . . , DEV(i), DEV(j), and DEV(k) is a memory controller and memory devices such as DRAMs and flash memories. FIGS. 1A and 1B illustrate a memory controller as the master device 100 and memory devices as the slave devices DEV(1), DEV(2), . . . , DEV(i), DEV(j), and DEV(k).\nReferring to FIG. 1B, the respective slave devices DEV(1), DEV(2), . . . , DEV(i), DEV(j), and DEV(k) are stacked and formed, and are coupled to the master device 100 through an interposer 110. Pillars formed through the stacked slave devices DEV(1), DEV(2), . . . , DEV(i), DEV(j), and DEV(k) include TSVs through which signals (data) are transmitted, that is, which form channels. The entire system illustrated in FIG. 1B may be implemented in one semiconductor chip package, and only the stacked slave devices may be implemented in one semiconductor chip package.\nFIG. 2 is a diagram illustrating the channels formed as the TSVs between the master device 100 and the slave devices DEV(i), DEV(j), and DEV(k) and RLC (resistance, inductance, and capacitance) components of the channels.\nReferring to FIG. 2, the respective channels have the RLC components. Therefore, the signals (data) transmitted through the channels may be delayed. The delay increases in proportional to the distance between the devices. That is, a flight time of a signal between the master device 100 and the slave device DEV(i) is longer than a flight time of a signal between the master device 100 and the slave device DEV(k). For reference, in FIG. 2, Tx represents transmission terminals provided in the master device 100 and the slave devices DEV(i), DEV(j), and DEV(k), and Rx represents reception terminals provided in the master device 100 and the slave devices DEV(i), DEV(j), and DEV(k).\nFIG. 3 is a diagram showing timing variations in signal transmission between the master device 100 and the slave devices DEV(i), DEV(j), and DEV(k), which occur depending on differences in channel length.\nIn FIG. 3, CMD represents a command applied from the master device 100 to the slave devices DEV(i), DEV(j), and DEV(k), D represents data which the master device 100 transfers to the slave devices DEV(i), DEV(j), and DEV(k), and Q represents data which the slave devices DEV(i), DEV(j), and DEV(k) transfer to the master device 100. The data Q are generated when the slave devices DEV(i), DEV(j), and DEV(k) processes the data D according to the command D.\nReferring to FIG. 3, the command and the data D transferred from the master device 100 to the slave device DEV(i) have a flight time of X(i). When the slave device DEV(i) transfers the data Q to the master device 100 in response to the command CMD and the data D, the data Q has a flight time of X(i). Therefore, while the data are exchanged between the master device 100 and the slave device DEV(i), the flight time is 2*X(i). Similarly, a flight time between the master device 100 and the slave device DEV(j) is 2*X(j), and a flight time between the master device 100 and the slave device DEV(k) is 2*X(k).\nThat is, the flight time of the signal may differ depending on with which device the master device 100 communicates among the slave devices DEV(i), DEV(j), and DEV(k), and the timing of transmission or reception of a specific signal (data) may be significantly varied."}
{"text": "Medical beds or tables used in hospitals, clinics, doctor offices or other medical environments may include a rail operatively secured to the bed for supporting various medical instrumentation. The rail may be used to secure a robotic catheter system as illustrated in U.S. Pat. No. 8,480,618 entitled Catheter System incorporated herein by reference in its entirety. A medical device mounted to the rail can apply a force on the rail that may cause torsional deflection of the interface between the rail and the frame of the bed or table.\nIt would be desirable to provide an apparatus for reacting moments on a rail of a medical bed and to provide support to the rail when the rail is used to support medical instruments."}
{"text": "The present disclosure relates generally to dimming controllers and dimming methods, and, more particularly, to dimming controllers suitable of receiving a dimming signal no matter it is a pulse-width-modulation (PWM) signal or a direct-current (DC) signal.\nLight emitting diode (LED), due to its characteristics in high power efficiency, compact product size, and long lifespan, has been widely adapted by lighting appliances and backlight modules. Until recently, most of cold cathode fluorescent lamps (CCFL) in the backlight modules of TV or computer display panels, for example, are replaced by LED modules.\nLED modules usually need dimming controllers to perform light dimming, so as to adjust the luminance of a display panel for example. There are two different methods in the art to dim the luminance of a LED module: PWM dimming and DC dimming. PWM dimming, also named digital dimming, employs a PWM or digital signal that jumps quickly back-and-forth between levels of “0” and “1” in logic to determine the duty cycle of a LED module, the ratio of the time when the LED module emits light to the cycle time of the PWM signal. For example, when the PWM signal is “1” in logic, the luminance of the LED module is in its maximum, and when the PWM signal is “0”, it is zero, not emitting light. In other words, PWM dimming makes a LED module blinking. In contrast, DC dimming, also known as analog dimming or resistive dimming, makes a LED module emitting light continuously while the luminance of the LED module corresponds to the voltage level of a DC or analog signal.\nFor having more market share, a dimming controller should accommodate a dimming signal no matter the dimming signal is of PWM or of DC, and provide appropriate luminance control."}
{"text": "This invention relates in general to enclosures for electrical components. In particular, this invention relates to an apparatus for pivotably mounting such an electrical enclosure on a wall or other support surface.\nA variety of enclosures are known in the art for protectively enclosing one or more electrical components, such as switches, circuit breakers, and the like. Such electrical enclosures are designed to protect the enclosed electrical components from unauthorized access and from external environmental conditions. A variety of standards have been promulgated for such electrical enclosures that specify their design, construction, materials, and resulting level of protection that can be expected therefrom. Some of such standards prohibit the formation of any openings through the walls of such electrical enclosures because penetration through any of these walls could reduce the level of protection that is provided thereby. Examples of such electrical enclosure standards include NEMA (National Electrical Manufacturers Association) standard types 4, 4×, and 12.\nElectrical enclosures may be used to contain main control components, operator interface controls, or other devices that may need to be located on or near machinery in an industrial facility. Such electrical enclosures are usually located in close proximity to the related machinery, but should not prevent access to the machinery.\nThere are various systems in use today that use a group of components to allow motion of the enclosure relative to a stationary support surface. However, these systems are usually intended for use with operator interface enclosures. Some machine tool manufacturers have devised their own means of pivoting an enclosure, but these means are usually designed for a specific application and require machining, welding, or a modified enclosure to facilitate the installation. Large hinges have been used to hang boxes where the wiring was installed in flexible conduits, but, over time, the conduits can become broken due to fatigue or abuse, thus compromising the integrity of the wiring, as well as posing a safety concern. Thus, it would be desirable to provide an apparatus for pivotably mounting such an electrical enclosure on or near machinery such that the electrical enclosure can be pivoted to access to the related machinery behind the enclosure without damage or disconnection of the components contained within the electrical enclosure."}
{"text": "1. Field of the Invention\nThe present invention relates to improvements made to lids for color cans, of the type described in FR-A-2 600 975.\n2. Discussion of the Prior Art\nIt is known from the prior art to provide a lid having a pouring spout for cans of base colors in agitator machines, which comprises:\nat least one member for fastening on a color can; PA1 a pouring spout, the edge of the opening of which is strictly planar; PA1 an operating lever pivoted at its bottom end to the top wall of the lid and movable between a position of rest and a pulled operative position; PA1 a likewise plane slide slidable on said edge in such a manner as to close it leaktightly, and being pulled for pouring purposes by means of the operating lever, wherein the operating lever is provided on its front face, at a distance from the pouring spout, with a stop member securing said operating lever in the position of rest and bearing leaktightly against a vent opening formed in the top wall of the lid, the slide being fastened to said operating lever by a lever pivoted on the one hand at the top middle part of the slide and on the other hand to said operating lever at a distance from its pivot axis, while spring means return said slide under pressure against the opening edge of the pouring spout and return said operating lever to its position of rest, in which the opening edge of the pouring spout is closed by the slide. PA1 at least one member for fastening on a color can; PA1 a pouring spout, the edge of the opening of which is strictly planar; PA1 an operating lever pivoted at its bottom end to the top wall of the lid and movable between a position of rest and a pulled operative position; PA1 a likewise planar slide slidable on said edge in such a manner as to close it leaktightly, and being pulled for pouring purposes by means of the operating lever, the latter having on its front face, at a distance from the pouring spout, a stop member securing said operating lever in the position of rest and bearing leaktightly against a vent opening formed in the top wall of the lid, the slide being fastened to said operating lever by a lever pivoted on the one hand at the top middle part of the slide and on the other hand to said operating lever at a distance from its pivot axis, while spring means return said slide under pressure against the opening edge of the pouring spout and return said operating lever to its position of rest, in which the opening edge of the pouring spout is closed by the slide, in which lid the vent opening is provided on the inner face of the lid with a channel which leads out at a point close to the edge of the lid and substantially diametrically opposite the pouring spout.\nThis arrangement already enables the operating lever, in the position of rest, to close the vent opening of the color can, which thus becomes completely isolated from the outside atmosphere, and, in the operating position, to open this opening, by simultaneously pulling the pouring spout slide, thus permitting a regular flow of paint when the latter is poured out.\nNevertheless, since the vent opening is situated to the rear of the agitator axis, relative to the pouring spout, and relatively close to the axis, it is not possible to prevent paint from also flowing out of the vent opening when a full can of paint is sharply tipped, so that the lid is soiled and the regularity of flow of the poured paint is reduced."}
{"text": "Field of the Invention\nThe present invention concerns an electric switch device comprising separable contacts respectively disposed on a fixed contact conductive part and a mobile contact conductive part which are connected to input and output terminals of the device, an arc extinguishing chamber and a shunt conductor adapted to receive, after it has jumped a dielectric gap, one end of the arc generated on separation of the contacts and displaced towards the arc extinguishing chamber by the usual blast arrangements or phenomena.\nDescription of the Prior Art\nDevices of this kind are described in patent FR-A-2 584 529, in which the shunt conductor shunts and consequently protects electromechanical and electronic switching units in series with the contacts concurrently with the arcing phenomenon.\nDocument EP-A-0 350 825 describes the combination of the shunt conductor, which in this instance includes the arc extinguishing device, with a current limiter device opening electrodynamically in the event of an overcurrent. Operation of the electrodynamically opening device is conditioned by the growth of the arc in the arc splitter plates which constitute the arc extinguishing device and therefore by the configuration and the performance of this device.\nAn object of the invention is, in an electromechanical switch device of the type described, to combine a reduction in the loading of the contacts by the arc and a limitation of the transient shunt current, independently of the arc extinguishing device."}
{"text": "This invention relates to double action presses used to draw hollow articles from a blank, and more particularly but not exclusively to long stroke double action mechanically driven presses used to redraw and wall iron a cup shaped blank, to make a can body.\nU.S. Pat. No. 3,696,657 (MAYTAG) describes a press having a first action to drive a ram and second action rod for driving a blank holder to co-operate with dies to redraw and wall iron a cup to a taller can body. The blank holder is mounted on a redraw carriage which is connected to the second action rod by pivotable joints which may wear to spoil the accuracy of blank holder motion. EP-A-0536812 describes a modified \"MAYTAG\" press in which the weight of the second action rod is supported to reduce wear on the pivotable joints and the redraw carriage is guided by pillars mounted in the die housing. The axis of each pivotable joint is arranged perpendicular to the driving rod to act in the manner of a universal joint permitting tilt but not lateral deviation of the rod direction from the redraw carriage direction by the pillars.\nOther longstroke mechanical presses for the purpose of drawing and wall ironing can bodies are described in U.S. Pat. Nos. 3,270,544 (Maeder), 3,663,072 (Cvacho) 3,735,629 (Paramonof) and 4,173,138 (Main). These mechanical presses have served the can making industry well and improvements have permitted increase of stroking speed from about 100 stokes per minute in 1970 to over 300 strokes per minute in 1990. However, there persists a need to increase press speed and reduce lost time arising from tool changes.\nIn these prior art presses a long ram is urged to move towards and away from a tool support, fixed to the press frame, by a connecting coupling pivotally connected to one end of a lever which swings on a pivot fixed to the press frame. The lever is driven to swing by a connecting rod attached to a crank shaft to constitute a first action linkage. When provided a second action, such as described in U.S. Pat. No. 3,735,629, comprises a pair of second levers pivoted on the frame and driven to swing by cams on the crank shaft to urge a cross head to move a pair of long connecting rods to move a second cross head towards and away from the tool support.\nA tubular blank holder is mounted on this second cross head to cooperate with a punch on the ram to permit redrawing of a shallow drawn cup against a redrawing die in the tool support. After redrawing to reduced overall diameter and increased height of side wall, the redrawn cup is pushed through at least one ironing ring to make a can body having a side wall thinner than its bottom wall. Typically the press stroke is about 20 inches long so that the first action ram is over 40\" long. Typically, the connecting rods of the second action are also about 40\" long. As press speeds increase the inertia forces arising from the combined weight of second action cross heads and connecting rods also increase so an improved support of the blank holder is required. Furthermore, adjustment of the position of the blank holder in respect to the redrawing die and ironing rings required laborious adjustment of the blank holder using adjusting the second crosshead, whenever the blank holder is changed, to ensure application of a correct blank holding force in alignment with the die and rings. Incorrect alignment of the blank holder with a redrawing die gives rise to an uneven holding pressure around the blank so that the height of the redrawn cup varies around the side wall of the redrawn cup."}
{"text": "Businesses use contact centers that include automated systems and representatives of the business to process transactions and/or service the needs of their customers. Contact centers may use a number of communication channels to engage the customers, such as telephone, email, live web chat, and the like.\nAn interaction from a customer such as, for example, an inbound voice call, is routed to a contact center target, such as, for example, a contact center agent, for processing. In simplest terms, routing is a process of sending the interaction to the target, such as, for example, sending an incoming telephone call or an incoming e-mail, to an agent. In practice, many steps are taken between the arrival of the interaction and the selection and use of a target. Generally speaking, not all interactions should go to the same target; choices should be made in order to determine the best target for each interaction. Each choice-point is an opportunity to make a decision based on the current situation—with the general goal of getting the interaction delivered to the right target."}
{"text": "It is generally desired to identify antibodies that specifically bind to particular classes of cells (e.g. tumor cells). Such antibodies, sometimes referred to as “targeting antibodies”, can be used to specifically direct (deliver) various effector molecules (e.g. liposomes, cytotoxins, labels, etc.) to the target cell.\nTargeting antibodies have been of particular interest in the study and treatment of cancer. A major goal of cancer research has been to identify tumor antigens that are qualitatively or quantitatively different from normal cells (Goldenberg (1994) Ca: A Cancer J. for Clinicians. 44: 43-64). The presence and/or quantity of such antigens could be detected by antibodies and such detection forms the basis of diagnostic and prognostic tests. In addition, the antibodies could be used to selectively kill tumor cells either directly via their effector function (Brown et al. (1989) Blood. 73: 651-661) or by attaching cytotoxic molecules to the antibody (Vitetta et al. (1987) Science. 238: 1098-1104; Brinkmann et al. (1993) Proc. Natl. Acad. Sci. USA. 90: 547-551).\nDespite the demonstration of antigens that are overexpressed on tumor cells, antibodies have been used with limited success for diagnosis and treatment of solid tumors, (see review in Riethmuller et al. (1992) Curr. Opin. Immunol. 4: 647-655, and Riethmuller (1993) Curr. Opinion Immunol. 5: 732-739). Their utility has been hampered by the paucity of tumor specific antibodies, antibody immunogenicity, low binding affinity, and poor tumor penetration.\nNonspecific toxicity results from the failure of the antibody to bind specifically and with high affinity to tumor cells. As a result, nonspecific cell killing occurs. In addition, the foreign immunotoxin molecule elicits a strong immune response in humans. The immunogenicity of the toxin portion of the immunotoxin has recently been overcome by using the human analogue of RNase (Rybak et al. (1992) Proc. Nat. Acad. Sci., USA, 89: 3165). The murine antibody portion, however, is still significantly immunogenic (Sawler et al., (1995) J. Immunol., 135: 1530).\nImmunogenicity could be avoided and toxicity reduced if high affinity tumor specific human antibodies were available. However, the production of human monoclonal antibodies using conventional hybridoma technology has proven extremely difficult (James et al., (1987) J. Immunol. Meth., 100: 5). Furthermore, the paucity of purified tumor-specific antigens makes it necessary to immunize with intact tumor cells or partially purified antigen. Most of the antibodies produced react with antigens that are also common to normal cells and are therefore unsuitable for use as tumor-specific targeting molecules."}
{"text": "The present invention relates to magnetic resonance imaging (MRI) systems, and particularly to the radio frequency (RF) coils used in such systems.\nMagnetic resonance imaging (MRI) utilizes hydrogen nuclear spins of the water molecules in the human body or other tissue, which are polarized by a strong, uniform, static magnetic field generated by a magnet (referred to as B0— the main magnetic field in MRI physics). The magnetically polarized nuclear spins generate magnetic moments in the human body. The magnetic moments point in the direction of the main magnetic field in a steady state, and produce no useful information if they are not disturbed by any excitation.\nThe generation of nuclear magnetic resonance (NMR) signal for MRI data acquisition is achieved by exciting the magnetic moments with a uniform radio frequency (RF) magnetic field (referred to as the B1 field or the excitation field). The B1 field is produced in the imaging region of interest by an RF transmit coil which is driven by a computer-controlled RF transmitter with a power amplifier. During the excitation, the nuclear spin system absorbs magnetic energy, and its magnetic moments precess around the direction of the main magnetic field. After the excitation, the precessing magnetic moments will go through a process of free induction decay, emitting their absorbed energy and then returning to the steady state. During the free induction decay, NMR signals are detected by the use of a receive RF coil, which is placed in the vicinity of the excited volume of the human body. The NMR signal is an induced electrical motive force (voltage), or current, in the receive RF coil that has been induced by the flux change over some time period due to the relaxation of precessing magnetic moments in the human tissue. This signal provides the contrast information of the image. The receive RF coil can be either the transmit coil itself, or an independent receive-only RF coil. The NMR signal is used for producing magnetic resonance images by using additional pulsed magnetic gradient fields, which are generated by gradient coils integrated inside the main magnet system. The gradient fields are used to spatially encode the signals and selectively excite a specific volume of the human body. There are usually three sets of gradient coils in a standard MRI system, which generate magnetic fields in the same direction of the main magnetic field, varying linearly in the imaging volume.\nIn MRI, it is desirable for the excitation and reception to be spatially uniform in the imaging volume for better image uniformity. In a standard MRI system, the best excitation field homogeneity is usually obtained by using a whole-body volume RF coil for transmission. The whole-body transmit coil is the largest RF coil in the system. A large coil, however, produces lower signal-to-noise ratio (S/N) if it is also used for reception, mainly because of its greater distance from the signal-generating tissues being imaged. Since a high signal-to-noise ratio is the most desirable factor in MRI, special-purpose coils are used for reception to enhance the S/N ratio from the volume of interest.\nIn practice, a well-designed specialty RF coil should have the following functional properties: high S/N ratio, good uniformity, high unloaded quality factor (O) of the resonance circuit, and high ratio of the unloaded to loaded Q factors. In addition, the coil device must be mechanically designed to facilitate patient handling and comfort, and to provide a protective barrier between the patient and the RF electronics. Another way to increase the S/N is by quadrature reception. In this method, NMR signals are detected in two orthogonal directions, which are in the transverse plane or perpendicular to the main magnetic field. The two signals are detected by two independent individual coils which cover the same volume of interest. With quadrature reception, the S/N can be increased by up to √2 over that of the individual linear coils.\nTo cover a large field-of-view, while maintaining the S/N characteristic of a small and conformal coil, a linear surface coil array technique was created to image the entire human spines (U.S. Pat. No. 4,825,162). Subsequently, other linear surface array coils were used for C.L. spine imaging, such as the technique described in U.S. Pat. No. 5,198,768. These two devices consist of an array of planar linear surface coil-elements. These coil systems do not work well for imaging deep tissues, such as the blood vessels in the lower abdomen, due to sensitivity drop-off away from the coil surface.\nTo image the lower extremities, quadrature phased array coils have been utilized such as described in U.S. Pat. Nos. 5,430,378 and 5,548,218. The first quadrature phased array coil, images the lower extremities by using two orthogonal linear coil arrays: six planar loop coil elements placed in the horizontal plane and underneath the patient and six planar loop coil elements placed in the vertical plane and in between the legs. Each linear coil array functions in a similar way as described in U.S. Pat. No. 4,825,162 (Roemer). The second quadrature phased array coil (Lu) was designed to image the blood vessels from the pelvis down. This device also consists of two orthogonal linear coil arrays extending in the head-to-toe direction: a planar array of loop coil elements laterally centrally located on top of the second array of butterfly coil elements. The loop coils are placed immediately underneath the patient and the butterfly coils are wrapped around the patient. Again, each linear coil array functions in a similar way as described in U.S. Pat. No. 4,825,162.\nIn MRI, gradient coils are routinely used to give phase-encoding information to a sample to be imaged. To obtain an image, it is required that all the data points in a so-called “k-space” (i.e., frequency space) must be collected. Recently, there have been developments where some of the data points in k-space are intentionally skipped and at the same time use the intrinsic sensitivity information of RF receive coils as the phase-encoding information for those skipped data points. This action takes place simultaneously, and thus is referred to as partially parallel imaging or partially parallel acquisition (PPA). By collecting multiple data points simultaneously, it requires less time to acquire the same amount of data, when compared with the conventional gradient-only phase-encoding approach. The time savings can be used to reduce total imaging time, in particular, for the applications in which cardiac or respiratory motions in tissues being imaged become concerns, or to collect more data to achieve better resolution or S/N. SiMultaneous Acquisition of Spatial Harmonics, SMASH, (U.S. Pat. No. 5,910,728 and “Simultaneous Acquisition of Spatial Harmonics (SMASH): Fast Imaging with Radiofrequency Coil Arrays,” Daniel K. Sodickson and Warren J. Manning, Magnetic Resonance in Medicine 38:591–603 (1997), both incorporated herein by reference) and “SENSE: Sensitivity Encoding for Fast MRI,” Klaas P. Pruessmann, et al., Magnetic Resonance in Medicine 42:952–962 (1999, also incorporated by reference, are basically two methods of PPA. SMASH takes advantage of the parallel imaging by skipping phase encode lines that yield decreasing the Field-of-View (FOV) in the phase-encoding direction and uses coils (e.g., coil arrays) together with reconstruction techniques to fill in the missing data points in k-space. SENSE, on the other hand, is a technique that utilizes a reduced FOV in the read direction, resulting an aliased image that is then unfolded in x-space (i.e., real space), while using the RF coil sensitivity information, to obtain a true corresponding image. Here, we make use of phase difference between signals from multiple coils to skip phase encoding steps. By skipping some of the phase encoding steps, one can achieve speeding up imaging process by a reduction factor R. Theoretically speaking, the factor R should equal the number of independent coils/arrays. In the SENSE approach, the SIR is defined as:                SNRSENSE=SNRFULL/{g√R}where SNRFULL is the S/N achievable when all the phase encoding steps are collected by traditional gradient phase encoding scheme. SNRSENSE is optimized when the geometry factor g equals 1. To obtain g of 1, traditional decoupling techniques such as overlapping nearest neighbor elements to null the mutual inductance between them shall not apply, as have been reported by others.        \nSENSE and SMASH or a hybrid approach of both demand a new type of design requirements in RF coil design. In SMASH, the primary criterion for the array is that it be capable of generating sinusoids whose wavelengths are on the order of the FOV. This is how the target FOV along the phase encoding direction for the array is determined. Conventional array designs can incorporate element and array dimensions that will give optimal S/N for the object of interest. In addition, users of conventional arrays are free to choose practically any FOV, as long as severe aliasing artifacts are not a problem. In contrast, when using SMASH, the size of the array determines the approximate range of FOVs that can be used in the imaging experiment. This then determines the approximate element dimensions, assuming complete coverage of the FOV is desired, as in most cases. In SENSE, the method is based upon the fact that the sensitivity of a RF receiver coil generally has a phase-encoding effect complementary to those achieved by linear field gradients. For SENSE imaging, the elements of a coil array should be smaller than for common phased-array imaging, resulting in a trade-off between basic noise and geometry factor, and adjacent coil elements should not overlap for a net gain in S/N due to the improved geometry factor when using SENSE.\nFor PPA applications, different types of RF coils or arrays have been used so far. However, most of them are based upon “traditional” RF coil design requirements, thus remain within the conventional coil design scheme. It has been reported, however, that since the phase information of B1 of a receive coil is very important when SENSE applications are demanded, for example, new coil design techniques such as non-overlapping adjacent coil elements may be necessary for better definition of the individual phase information associated with each RF coil used in an array, unlike traditional design scheme where two adjacent coils elements are overlapped to null the mutual inductance between the elements (U.S. Pat. No. 4,825,162). Without overlap, the coupling may be increased, but there is a net gain in S/N due to the improved geometry factor when using SENSE. As stated in the above, the use of smaller coil-elements than those for conventional imaging results in a trade-off between basic noise and geometry factor."}
{"text": "Voice-enabled systems help users complete assigned tasks. For example, in a workflow process, a voice-enabled system may guide users through a particular task. The task may be at least a portion of the workflow process comprising at least one workflow stage. As a user completes his/her assigned tasks, a bi-directional dialog or communication stream of information is provided over a wireless network between the user wearing a mobile computing device (herein, “mobile device”) and a central computer system that is directing multiple users and verifying completion of their tasks. To direct the user's actions, information received by the mobile device from the central computer system is translated into speech or voice instructions for the corresponding user. To receive the voice instructions and transmit information, the user wears a communications headset (also referred to herein as a “headset assembly” or simply a “headset”) communicatively coupled to the mobile device.\nThe user may be prompted for a verbal response during completion of the task. The verbal response may be a string of characters, such as digits and/or letters. The string of characters may correspond, for example, to a credit card number, a telephone number, a serial number, a vehicle identification number, or the like. The spoken string of characters (i.e., the verbal response) may be referred to herein as a “spoken complete value entry”.\nUnfortunately, speaking the spoken complete value entry is often time-consuming and difficult to speak correctly, especially if the spoken complete value entry is long (i.e., the number of characters in the spoken string of characters is relatively large (referred to herein as a “spoken long value entry”). If errors in speaking the spoken complete value entry are made by the user (i.e., the speaker), conventional systems and methods often require that the user restart the spoken string, causing user frustration and being even more time-consuming.\nTherefore, a need exists for auto-complete methods for spoken complete value entries, particularly spoken long value entries, and for use in a workflow process."}
{"text": "1. Field of the Invention\nThe present invention relates to a microcomputer and an electronic timepiece utilizing a microcomputer. More specifically, the present invention relates to a microcomputer employing an improved addressing system of a random access memory and an electronic timepiece utilizing such a microcomputer.\n2. Description of the Prior Art\nA conventional electronic timepiece employs frequency division utilizing counters. However, such conventional electronic timepiece involves an inconvenience in which the electronic circuit need be redesigned from the beginning in making modifications of the particulars of a timepiece circuit. In order to eliminate such inconvenience, an electronic timepiece utilizing a microcomputer has also been proposed and put into practical use. Such electronic timepiece utilizing a microcomputer is referred to as a microcomputer type timepiece in the present application for simplicity of description. Generally, a microcomputer type timepiece has an advantage that modifications of the particulars can be made by simply changing masks for use in fabricating a read only memory for storing a program.\nFIG. 1 is a block diagram of a conventional microcomputer type timepiece. The timepiece shown comprises a reference signal generating circuit 1, a clock generator 2, a switch input circuit 3, a state control circuit 4, a timer counter 5, a chronograph counter 6, a read only memory 7, a program counter 8, a stack 9, an instruction register 10, a control circuit 11, an arithmetic logic unit 12, an accumulator 13, a random access memory 14, a decoder 15, and a latch circuit 16.\nThe reference signal generating circuit 1 is connected to a quartz resonator 17 and makes oscillation to provide an oscillating signal the frequency of which is determined by the quartz resonator 17. The oscillating signal is frequency divided and a predetermined frequency divided signal is applied to the clock generator 2, the switch input circuit 3, the timer counter 5, the chronograph counter 6 and the like. The reference frequency generating circuit 1 is also adapted to provide an interrupt signal INT2 for the purpose of causing a time count operation at each predetermined period of time, say for every 0.5 second.\nThe clock generator 2 generates a system clock signal for performing an operation as a central processing unit by the use of a frequency divided signal obtained from the reference frequency generating circuit 1 and the system clock signal is applied to various portions of the circuit. The clock generator 2 is supplied with a stop signal CLKSTOP for stopping generation of the clock signal from the state control circuit 4, so that the clock generator 2 is responsive to the stop signal CLKSTOP to stop the operation. The clock generator 2 is also connected to receive from the reference signal generating circuit 1, the chronograph counter 6 and the like the interrupt signals INT0, INT1, INT2 and INT3 and the timer output TIMEROUT and to receive an operation start signal CLKSTART, so that the clock generator is responsive to these signals to restart the operation which has been stopped.\nThe switch input circuit 3 comprises external terminals M1 to M4 and S1 to S4 for preventing a chattering phenomenon of the switches connected thereto and for receiving the on/off data of the switches connected thereto, while the data is transferred to the data bus DB. The on/off control of the switches connected to the external terminals S1 to S4 can be designated by the switch control signal SWCON of the state control circuit 4 and can provide an operation start signal CLKSTART in response to the on/off of the designated switch. Furthermore an interrupt signal INT3 can also be obtained in response to the on/off of the switches connected to the external terminals S1 to S4. In the chronograph mode a chronograph control signal CHRCON for controlling start, lap and stop of the chronograph counter 6 is obtained in response to the on/off of the switches connected to the terminals S1 and S2.\nThe state control circuit 4 is controlled in accordance with the program obtained over the instruction bus IB so that the states of the various circuits are controlled in various modes and provides the switch control signal SWCON for controlling the switch input circuit 3, the timer control signal TIMCON for performing selection of the reset, start and timer time period of the timer counter 5 and the stop signal CLKSTOP for stopping the operation of the clock generator 2. Furthermore there are provided an external terminal LAMP for connection of a lamp and an external terminal ALM connected to an alarm tone generator.\nThe timer counter 5 can allow for setting of a timer time period in accordance with the program sent over the instruction bus IB and provides a signal TIMEROUT for starting the operation of the clock generator 2 after the lapse of the timer time period.\nThe chronograph counter 6 comprises a counter for counting up to 1/100 second and 1/10 second and provides an interrupt signal INT1 at every 1/10 second. The count value in the chronograph counter 6 is applied through the chronograph data bus CDB to the decoder 15.\nThe read only memory 7 comprises 1390.times.14 bits and stores in a fixed manner a program for performing a timing operation and various controls. To that end the read only memory 7 provides to the instruction register 10 an instruction code of 14 bits stored in the address designated by the program counter 8 of 11 bits.\nThe instruction register 10 stores an instruction code obtained from the read only memory 7 and provides the same on the instruction bus IB.\nThe program counter 8 receives the next address or jump control signal after processing of one instruction to provide to the read only memory 7 an address signal of the address to which jump is to be made.\nThe stack 9 has the capacity of eight levels and stores the sum of the count value in the program counter 8 plus one in the case where the operation should shift to the interrupt processing or subroutine. However, an interrupt processing is to be made after the instruction \"SCP\" for stopping the operation of the clock generator 2 is executed, the count value of the program counter 8 is as such stored in the stack 9 and, after the end of the interrupt processing, again the stop instruction \"SCP\" of the clock is addressed and the same is executed.\nThe random access memory 14 stores in appropriate regions the data concerning the current time, an alarm time, a timer time period and the like. The addresses of the random access memory 14 are designated by the instructions obtained from the instruction register 10 and the data stored in the designated addresses is sent to the data bus DB in response to the read/write control signal R/W obtained from the control circuit 11 or the data sent to the data bus DB is stored in the designated addresses.\nThe arithmetic logic unit 12 receives data sent onto the data bus DB and makes an arithmetic operation such as addition, subtraction and the like and the result is stored in the accumulator 13 and is sent again onto the data bus DB.\nThe decoder 15 and the latch circuit 16 constitute a display circuit. The decoder 15 converts the applied data of four bits to seven segment signals for display and provides the same to the segment bus SB. The latch circuit 16 holds the segment signals sent to the segment bus SB and provides the same to the external terminals a to g and to that end the external terminals of the number corresponding to that of the segments are provided. Which latch circuit should store the converted segment signals is determined in accordance with the program and, in the case where the data of the first digit of the current time is converted, for example, the latch circuit 16 corresponding to the display unit for displaying the first digit of the current time is designated by the program. On the other hand, the decoder 15 is also connected to the chronograph data bus CDB for supply of the output from the chronograph counter 6 as well as the data bus DB and the decoder 15 is selectively connected to either the data bus DB or the chronograph data bus in accordance with the program.\nThe control circuit 11 deciphers the instructions of the read only memory 7 sent over the instruction bus IB and controls the respective circuits as a function of the system clock obtained from the clock generator 2 for the purpose of performing the operation corresponding to the instructions. Furthermore, the program counter 8 and the read only memory 7 are controlled to execute the processing corresponding to the interrupt when the interrupt signals INT0, INT1, INT2 and INT3 are provided.\nNow a principal operation of the FIG. 1 timepiece will be described. When a power supply is turned on, an initial clear is rendered effective and the respective circuits are reset, while the program counter 8 designates the address zero of the read only memory 7. The address zero of the read only memory 7 stores the program for initial setting and therefore \"0\" or predetermined data is stored in the random access memory 14 through execution of the program, whereby initial setting is performed. Then the program of mode change is executed. This program serves to determine which mode, the current time mode, the alarm mode, the timer mode or the chronograph mode, or if another function has been employed, an additional mode, has been selected in accordance with the contents of the data representing the on/off of the switch applied to the data bus DB from the switch input circuit 3.\nIn the case where the current time mode is selected, for example, the current time is displayed. More specifically, the minute digit, the ten-minute digit, the hour digit, and the ten-hour digit stored in the predetermined address of the random access memory 14 are in succession addressed in accordance with the display instructions of the addressable random access memory 14 and the data is sent to the data bus DB. The segment signals as converted by the decoder 15 are stored in the latch circuit 16. The decoder 15 is in advance controlled by the instructions so that the data bus DB may be connected to the decoder 15. The latch circuit 16 serves to keep providing the stored segment signals to the external terminals a to g until the new segment data is renewed.\nThen a coincidence detection program for detecting coincidence of the data concerning the current time stored in the predetermined address of the random access memory 14 and the data concerning the alarm time is executed in the case where the alarm mode is selected, for example. Coincidence detection is performed in succession for the ten-hour digit, the hour digit, the ten-minute digit and the minute digit and, upon detection of coincidence in all digits, the alarm signal is obtained at the output terminal ALM of the state control circuit 4.\nMeanwhile, some type of electronic timepiece is adapted to have a so-called multialarm function. By multialarm function is meant a function capable of setting a plurality of alarm times per day. For example, a plurality of alarm times such as in the morning, day and evening may be set with such electronic timepiece and thus a multialarm function is of much convenience.\nA typical conventional microcomputer type timepiece having a multialarm function is adapted such that a plurality of alarm times are stored in different addresses of the random access memory and coincidence of the same with the current time is detected for each of the stored data. Accordingly, a coincidence detection program for detecting coincidence of the current time and the alarm times requires two or more programs of different addressing data, with the resultant disadvantage that the program becomes complicated and lengthy. However, coincidence detection of the current time with alarm times is common to any of the alarm times. Therefore, if an addressing scheme for different alarm times could be simplified, the coincidence detection program could be accordingly simplified."}
{"text": "Biofuels have a long history ranging back to the beginning on the 20th century. As early as 1900, Rudolf Diesel demonstrated at the World Exhibition in Paris, France, an engine running on peanut oil. Soon thereafter, Henry Ford demonstrated his Model T running on ethanol derived from corn. Petroleum-derived fuels displaced biofuels in the 1930s and 1940s due to increased supply and efficiency at a lower cost.\nMarket fluctuations in the 1970s, due the Arab oil embargo and the Iranian revolution, coupled to the decrease in U.S. oil production led to an increase in crude oil prices and a renewed interest in biofuels. Today, many interest groups, including policy makers, industry planners, aware citizens, and the financial community, are interested in substituting petroleum-derived fuels with biomass-derived biofuels. The leading motivation for developing biofuels is of economical nature, namely, the threat of ‘peak oil’, the point at which the consumption rate of crude oil exceeds the supply rate, thus leading to significantly increased fuel cost, and resulting in an increased demand for alternative fuels.\nBiofuels tend to be produced with local agricultural resources in many, relatively small facilities, and are seen as a stable and secure supply of fuels independent of geopolitical problems associated with petroleum. At the same time, biofuels enhance the agricultural sector of national economies. In addition, environmental concerns relating to the possibility of carbon dioxide related climate change is an important social and ethical driving force which is starting to result in government regulations and policies such as caps on carbon dioxide emissions from automobiles, taxes on carbon dioxide emissions, and tax incentives for the use of biofuels.\nThe acceptance of biofuels depends primarily on economical competitiveness of biofuels when compared to petroleum-derived fuels. As long as biofuels cannot compete in cost with petroleum-derived fuels, use of biofuels will be limited to specialty applications and niche markets. Today, the use of biofuels is limited to ethanol and biodiesel. Currently, ethanol is made by fermentation from corn in the US and from sugar cane in Brazil and is competitive with petroleum-derived gasoline, exclusive of subsidies or tax benefits, if crude oil costs above $50 USD per barrel and $40 USD per barrel, respectively. Biodiesel has a breakeven price of crude oil of over $60 USD/barrel to be competitive with petroleum-based diesel (Nexant Chem Systems. 2006. Final Report, Liquid Biofuels: Substituting for Petroleum, White Plains, N.Y.)."}
{"text": "There are many industrial processes and applications that require a continuous supply of a process fluid used for cooling purposes. The devices used to furnish this process fluid supply are ordinarily called \"water chillers\", even though the process fluid is most usually a mixture of water and other liquid, usually ethylene glycol, and that term is used in this specification to refer to controlled chillers using any appropriate coolant.\nA water chiller must maintain the temperature of the process fluid essentially constant within a very limited tolerance range; in some instances the critical parameter may be selected as the return temperature from the process apparatus, but more frequently the supply temperature should be selected as the basis of control. The water chiller is often a self-contained unit that can be transported from one location to another, literally constituting a heat transfer system on wheels. Water chillers of this kind are most frequently employed in processes involving heat transfer where the coolant must be maintained constant at some temperature between -30.degree. F and +60.degree. F.\nFor the most part, the controls for water chillers have constituted conventional thermostats, freezestats, and pressure-actuated switches connected directly in the electrical actuating circuits of the pumps, compressors, and other components of the water chilling apparatus. A further control that is often used is a float switch, connected to a coolant reservoir, employed to interrupt operation whenever the water supply is inadequate. Another control frequently used is a hot-gas bypass valve, used to control cooling capacity during intermittent or partial load conditions.\nFor many applications, these direct acting controls are adequate. However, they do not afford the precision control essential to some sensitive manufacturing and other industrial operations. Furthermore, controls of this nature do not provide comprehensive protection for the water chiller and the process apparatus it serves relative to possible malfunctions in the chiller itself or in the process apparatus. Moreover, the inherent inertia of these conventional controls makes it difficult to control chiller operation on the basis of coolant supply temperature; it is usually necessary to use the coolant return temperature as the control parameter."}
{"text": "The present invention relates to apparatus and methods for controlling the gas flow between two chambers in a mass spectrometer. According to an embodiment one or both of the chambers may comprise a vacuum chamber.\nMass spectrometers often contain different regions or chambers which are at different levels of vacuum. For example, a mass spectrometer may comprise a quadrupole mass filter (“QMF”) which resides in a chamber at a pressure of approx, 1×10−5 mbar and which is followed by a collision cell at a pressure of approx. 1×10−3 to approx. 1×10−2 mbar. This in turn may be followed by a Time of Flight (“TOF”) mass analyser which may be operated at a pressure of <1×10−6 mbar.\nBetween these different regions there is normally an opening or differential pumping aperture which acts to limit the flow of gas from one chamber to another and through which ions must pass if they are to traverse the mass spectrometer. These openings are generally manufactured to be as thin as possible, typically 0.5 mm to 1.0 mm, so as to minimise loss of ion transmission as ions pass through the orifice. The thicker the opening is the more likely it is that some ions will strike the inner wall of the opening as they pass through the orifice and be lost.\nReducing the size of an opening (i.e. the diameter of a circular hole or the length of a slit) reduces the gas flow through it, which in turn reduces the quantity of vacuum pumping that is required to maintain the desired pressure in the different regions. This is particularly important in situations where there is a large pressure differential between vacuum chambers and hence a large gas flow, or where a small, lightweight or portable instrument is desired. However, reducing the size of an orifice makes it more difficult to focus ions through it. This can lead to ions no longer being able to pass through the orifice which in turn reduces the transmission and hence sensitivity of the mass spectrometer.\nIt is known to use a valve to reduce the gas flow into the initial vacuum chamber of a mass spectrometer from the atmosphere.\nIt is desired to provide an improved mass spectrometer and method of mass spectrometry."}
{"text": "It is known to use free electron laser (FEL) radiation sources to produce radiation of a desired wavelength, in which an electron beam comprising a periodic sequence of electron bunches is passed through an undulator to generate the radiation. Such sources can be used to produce radiation in a range 4 nm to 25 nm, for example extreme ultra-violet (EUV) radiation.\nIn known FEL radiation sources, ions are produced from residual gas in the electron beam through collisional ionization. An expected rate of ion generation in EUV FEL sources with a beam energy of a few hundred MeV and a beam current of a few tens of mA is such that, in the absence of counter-measures, the electron beam could be fully neutralized (e.g. ion charge matches electron charge per meter of electron beam) in less than 10 seconds. As a result, emittance of electron bunches at the undulator could significantly exceed 10 mm mRad and conversion efficiency would be expected to drop more than ten times, thus effectively disabling the operation of the EUV FEL.\nIn order to counteract ion concentration build-up several strategies have been suggested, including placing extraction electrodes along the electron beam or implementing a beam current pattern with short additional gaps between selected pairs of electron bunches (also referred to as clearing gaps). The short additional gaps can also be thought of as representing missing bunches in the sequence of electron bunches. The use of clearing gaps is intended to give ions, otherwise trapped in the potential of the bunched electron beam, some time to drift away. The use of clearing gaps according to known suggestions may not completely remove trapped ions but may allow mitigation of fast ion instability which can be detrimental for FEL operation.\nKnown FEL sources include LINACs for accelerating (and decelerating) electron bunches before (and after) they pass through the undulator. Energy recovery LINACs can be used, which are usually designed to operate with a balanced cavity load close to zero (e.g. currents in accelerating and decelerating beams match, and energy extracted and deposited upon acceleration and deceleration almost match).\nThe use of clearing gaps between electron bunches may distort operation of energy recovery LINACs. In the case of a single-pass LINAC system it has been suggested in G. H. Hoffstaetter et al, Nuclear Instruments & Methods in Physics Research A, 557 (2006), 205-212 to match clearing gaps (e.g. trains of missing bunches) in accelerating and decelerating electron beams within LINAC modules in a single pass LINAC. The matching of clearing gaps in a single pass LINAC may be relatively straightforward, given that each electron bunch circulates once around a loop and only passes through the single LINAC once when it is accelerating and once when it is decelerating, and may be achieved by a requirement that missing bunches are provided at a suitable regular rate.\nIt is an aim of the present invention to provide an improved, or at least alternative, FEL radiation source and method of operation of such a source that can provide for reduction of ion build-up."}
{"text": "This invention relates to the application of colorant in a desired pattern to an embossed electrically uncharged polymeric surface.\nIt is often desirable to decorate items having polymeric surfaces, such as suitcases, automobile seats, and furniture upholstery with colorants. Where a surface is smooth, these are various methods to decorate it including spraying with spray masks, lamination, printing, and the like. However, if the polymeric surface is embossed, lamination, spraying with spray masks, and printing methods are not highly desirable. Lamination and spray with spray masks will not seal against the various peaks and valleys of an embossed surface and printing methods also have disadvantages. Electrostatic printing methods normally require that the substrate, in this case the polymeric surface, be electrostatically neutral or that a uniform field can be produced between the printing screen and the substrate surface. To obtain a uniform electrical field between an embossed surface and the printing screen is difficult because of the uneven surface. Moreover, when a non-uniform electrostatic field is established between the substrate and the printing screen, the uneven surface may cause the designs to have feathered edges. Thus, designs are not as attractive and distinct as desired.\nThe prior art, e.g., U.S. Pat. No. 3,058,444 uses a rotatably mounted magnetic roller to deposit a mixture of magnetic powder and toner to an electrophotographically formed electrostatic latent image. The necessity of charging the surface and application of electrophotographic coating to the surface requires specialized equipment with its appurtenant cost. Moreover, variations in the amount of the electrostatic charge in the latent image could vary the concentration of the deposited mixture and/or cause a feathered effect at the edge of a design. Coatings have been applied to embossed surfaces but generally an adhesive method is used, e.g., see U.S. Pat. No. 3,616,135.\nThe present invention has the distinct advantage of not requiring the substrate to be electrically charged. That is, the material which has a polymeric surface is electrically uncharged although alternately it may be charged if desired. The theory on which the present invention is based is that if a colorant made of a polymeric material is pressed against another polymeric surface, the friction itself will cause an electrostatic attraction known as a triboelectric effect. This triboelectric effect at least temporarily holds the colorant against the polymeric surface until it can be subsequently coated with a coating of some type such as hydroxy functional acrylic with urethane and 1,6 diacrylate.\nAnother apsect of this invention is the use of a roller magnet bounded by generally non-magnetic walls on both sides thereof which extend radially outwardly so that the magnetic carrier and polymeric colorant is attracted and distributed in a straight line. The wheel-shaped magnet is biased against the polymeric surface so that it presses the mixture against the surface and deposits it thereon. The walls prevent any feathering of the stripe or other desired design.\nDifferent methods for holding a polymeric material in position on the polymeric surface are the applications of heat and resins. However, if heat and resins are not used in precise amounts, they will distort the polymeric surface. The subject invention has the advantage of less risk of damage to the surface. The coating is simply applied and permitted to dry by dehydration or by application of ultraviolet light."}
{"text": "An engine cooling fan and electric motor assembly is used to provide radiator and A/C condenser cooling airflow in engine compartments of many vehicles. The engine compartment itself is a high-temperature environment. To complicate matters, the motor generates its own internal heat during the normal course of operation. The internally generated heat is primarily due to Joule heating caused by electric current passing through its rotor and stator coils, and brush losses. High temperatures can greatly increase the wear of motor brushes and bearings, and significantly reduce motor life.\nOne solution to these problems is to provide cooling holes 44a, 44b in FIG. 4, to admit cooling air through the motor 22. The invention provides additional strategies to cool motor 22."}
{"text": "The present invention relates to a control device which can be fitted to a tap for regulating and cutting off a flow of combustible gas, particularly for domestic cookers.\nMore precisely, the invention relates to a control device of the type including a base on which a knob, which is intended to be fixed to a member for operating the tap, is rotatable to regulate the gas-flow."}
{"text": "1. Field of the Invention\nThe invention generally relates to equipment in baseball and softball. More specifically, the invention relates to bases that are used in the games of baseball and softball.\n2. Discussion of Related Art\nIn the United States and Canada and in many East Asian, Latin American and Caribbean countries, the games of softball and baseball are very popular.\nFamiliarity with the basic concepts of the rules of baseball and softball are assumed, as one of ordinary skill in the art would know the rules intimately. However, the rules of baseball and softball provide for markers known as bases that are positioned at the four corners of a diamond, up to 90 feet in distance depending upon the age level of the players or whether baseball of softball is being played. Three of the bases are commonly referred to as bags, because they are short boxes that are placed on the surface of the ground and rise up one to three inches in height. The fourth base, known as home plate, is much thinner and rises above the ground less than ½ inch.\nThese bases are commonly held in position in a number of ways. For example, major league teams and other organized baseball leagues generally use hollow, square metal posts or anchors located in the ground at the proper base locations. Stanchions that are attached to the bottom of the bases are placed in the hollow posts/anchors. In this manner, the bases are firmly and fixedly, yet removably positioned for play. Runners naming to the bases do not have to fear that the bases will slide as the runner slides into or steps on the base. This method requires digging a hole into the ground, placing a concrete bed in the hole, setting the anchor/post in the concrete bed and repacking the dirt around the post/anchor. This can be accomplished on a dedicated baseball field.\nA second common type of base uses spikes that are driven into the ground that are attached to straps. The straps are cinched tight to the base and the spikes that hold the base in position.\nA third common category of common bases is called “throwdown bases.” These bases typically are formed from a single piece or molded rubber or vinyl. These bases are generally thinner than, official bases.\nAnother type of base uses a magnet to secure a base having a bottom metal material. Dislodgement of the base exposes the magnet which can be poorly seen due to dirt and dust accumulating on it.\nOther types of bases use either an oversized hole to receive an anchor for the base or can use protruding rubber tips which are inserted in openings formed in a substrate and trapped therein until a sufficient force is applied to the base. Such rubber tips are not able to be freely removed from the substrate.\nMoreover, youth baseball leagues have recently mandated that a base is displaced from its position when a runner slides into it. Put another way, leagues are requiring a base that will move when a lateral load is applied. The difficulty is designing such a base that will not substantially move when the runner merely runs across the base, i.e. when the load applied is a combination of a lateral load and a vertical load.\nTherefore, there is a need in the art for detachable/removable baseball and softball bases that will not substantially move when a runner runs across it and yet will detach when a runner slides in with sufficient force."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device of a chip-on-chip structure that has a semiconductor chip bonded with other semiconductor chips thereon, and to a method for manufacturing same.\n2. Description of Related Art\nThere are semiconductor devices increased in integration degree, including system-on-chips (SOCs) and multi-chip modules (MCMs). In the system-on-chip, the functions conventionally realized on a plurality of ICs are integrated on one semiconductor chip. Meanwhile, the multi-chip module is structured by a plurality of semiconductor chips arranged with density on a wiring board of glass-epoxy or the like. Each of them has multiple functions as one semiconductor device and can be size-reduced as compared to a combination of a plurality of semiconductor devices realizing the equivalent functions. Meanwhile, this reduces the wiring length in the overall, enabling high-speed transmission of signals.\nHowever, the system-on-chip is complicated in manufacture process, requiring a huge amount of capital investment and hence high manufacture cost. Meanwhile, the multi-chip-module has a plurality of semiconductor chips mutually arranged side by side on a wiring board. Because these semiconductor chips are connected by wirings, the size is greater as compared to the system-on-chip and hence integration degree lowers.\nOn the other hand, there is a chip-on-chip structure as one form of an integration increased semiconductor device. The chip-on-chip-structured semiconductor device has a structure having a plurality of semiconductor chips that are oppositely placed for mutual connection. Such a semiconductor device is not configured with an integration of the functions conventionally realized on a plurality of ICs (semiconductor chips) as in the system-on-chip. Consequently, the manufacture process is not so complicated as that of the system-on-chip. Therefore, there is a merit to reduce the manufacturing cost.\nThe chip-on-chip-structured semiconductor devices include those having a plurality of small semiconductor chips (child chips) arranged side by side on one large semiconductor chip (parent chip). Such a semiconductor device, at a glance, is similar in structure to the multi-chip module having a plurality of semiconductor chips arranged side by side on a wiring board.\nIn the chip-on-chip-structured semiconductor device, however, the parent chip not only serves as a wiring board to connect between the child chips but also acts itself as a semiconductor chip having functional elements. This results in higher integration degree. On the other hand, the wiring formed on the parent chip, made by a semiconductor process, is by far finer than the wiring of a wiring board of a multi-chip-module. This allows the functional elements of the semiconductor chips (parent and child chips) to be connected at a short length of wiring, whereby signal transmission speed can be increased higher as compared to that of the multi-chip module.\nThe chip-on-chip-structured semiconductor devices include those that a child chip has, further, one or a plurality of child chips vertically superposed thereon. Namely, such a semiconductor device has a structure connected, on a parent chip, with one or a plurality of chip blocks each having one or a plurality of superposed semiconductor chips. This structure realizes a highly integrated semiconductor device.\nIn such a semiconductor device, however, the wiring between arbitrary two of the semiconductor chips is provided necessarily through a wiring plane (usually, active surface) of the parent chip, thereby increasing the mean wiring length. Namely, in the case the semiconductor chip is in a position close to the top end of the chip block (position distant from the parent chip), there is an increased length of wiring between that semiconductor chip and the parent chip. Consequently, signals could not be transmitted at sufficient high speed. Meanwhile, even if the wiring length is tried to decrease throughout the entire semiconductor device, there is encountered a low freedom in design, e.g. restriction in semiconductor chip arrangement."}
{"text": "The wide adoption of computer networking capability and processing power in industrial environments has resulted in the increased usage of industrial automation systems. These industrial automation systems can be complex and are expected to incorporate functional cooperation between different components of the industrial automation systems. The different components might be supplied by one vendor, or in many cases, more than one vendor. Therefore, there is often limited interoperability and consistency between components. On the user side, the limited interoperability and consistency has been blamed for user dissatisfaction. On the vendor side, interoperability with components from a different vendor is not necessarily a priority."}
{"text": "1. Technical Field\nThe present invention relates to a purge gas spraying plate and a fume removing apparatus having the same, and more particularly, to a purge gas spraying plate which effectively removes fumes by spraying a purge gas and a fume removing apparatus having the same.\n2. Description of the Related Art\nGenerally, semiconductor devices are manufactured by depositing thin films for performing various functions on upper surfaces of the wafer, and patterning the thin film to form various circuit geometries.\nUnit processes for manufacturing the semiconductor devices can be mainly divided into an impurity ion implantation process by which impurities such as group 3B elements or group 5B elements are implanted into an inside of the semiconductor, a thin film deposition process which forms material films on a semiconductor substrate, an etching process which forms a predetermined patterns using the material films, a planarization process such as a chemical mechanical polishing (CMP) process which removes step differences by polishing a surface of a wafer in a lump after depositing interlayer insulating films and the like on an upper surface of the wafer, and a cleaning process for removing residues.\nMeanwhile, wafers on which the predetermined process described above is completed are not transferred as an individual wafer, but are transferred as a batch of 20 to 25 wafers, which are seated in a cassette. At this time, the wafers on which process gases used for the process and fumes which are by-products of the process are not removed and remain on surfaces of the wafers are seated in the cassette.\nHowever, the residues described above causes the contamination of a semiconductor manufacturing apparatus and defects of etch patterns when a subsequent process is performed in a state in which the residues are attached to the surfaces of the wafers. As a result, since the reliability of a device is decreased, many techniques have been developed to remove the residues.\nAs described above, a conventional fume removing apparatus for removing fumes remaining on surfaces of wafers is disclosed in Korean Patent Application No. 1444241.\nAs illustrated in FIGS. 1 and 2, a cleaning apparatus 100 as a conventional fume removing apparatus is formed with a side cleaning apparatus 110 and a rear cleaning apparatus 130.\nThe side cleaning apparatus 110 includes wafer support members 111 and spacers 113.\nThe wafer support members 111 are each provided with a rib 111a and through holes 111b. \nThe rib 111a includes a step 111aa at an edge thereof to easily seat a wafer 5 inserted from one side thereof.\nThe plurality of through holes 111b are formed at equidistant intervals in a longitudinal direction to be easily coupled with a plurality of spacers 113, which will be described below using bolts BT.\nThe spacer 113 is provided on an upper side of the wafer support member 111.\nThe spacer 113 formed in a bar shape having a cross-sectional rectangular shape is provided with an inlet 113a in a longitudinal direction of one side thereof so that a purge gas sprayed from a pneumatic block flows in, and a plurality of spraying holes 113b at equidistant intervals in a direction perpendicular to a side of the inlet 113a. \nThe inlet 113a described above is directly connected to an outlet of the pneumatic block to enable a purge gas (N2 gas) to be individually sprayed through only the spraying holes 113b corresponding to a portion at which a wafer 5 is provided.\nThrough holes 113c are formed at locations respectively the same as the centers of the through holes 111b of the wafer support member 111 at equidistant intervals in a longitudinal direction of the spacer 113 to be easily coupled therewith using bolts.\nAs described above, a configuration of one set having the wafer support member 111 and the spacer 113 is vertically repeated as a plurality, and thus, it is possible to clean a plurality of stacked wafers 5.\nIn addition, the rear cleaning apparatus 130 is provided at one end of the side cleaning apparatus 110 in a rear direction, i.e., based on a direction illustrated in FIG. 1.\nThe rear cleaning apparatus 130 has a plurality of spraying holes 131, which are stacked and provided at a rear of the side cleaning apparatus 110 in a direction perpendicular thereto and formed at the same levels as the spraying holes 113b of the spacer 113, at equidistant intervals in a longitudinal direction.\nHowever, since a conventional fume removing apparatus is manufactured by separately manufacturing a wafer support member 111 which supports a wafer and a spacer 113, and installing a plurality of configurations of which one set includes the wafer support member 111 and the spacer 113 sequentially stacked at an inner side of the fume removing apparatus, there are problems in that a structure of the inside thereof is complex and the cost of manufacturing equipment is increased.\nIn addition, since the spraying holes 113b of the spacer 113 are formed to be only disposed linearly without considering a circular shape of a wafer, there is an area having a big gap between the spraying holes 113b and the circumference of the wafer, and thus, a dead zone in which fumes remaining on the wafer cannot be removed may be generated.\nIn addition, since a wafer is seated on a step 111aa of the rib 111a of the wafer support member 111, there is a problem in that a contact area in which the wafer and the step 111aa are in contact with each other is increased which can cause damage to the wafer."}
{"text": "In mineral extraction, particularly in underground mining, scraper chain conveyors, which can be configured as single scraper chain or double scraper chain conveyors, are used. Scraper chain conveyors of the generic type are used, in particular, as face conveyors, in which, over large face lengths, the mined material has to be continuously evacuated with the scrapers fastened to the revolving scraper chain and transferred onto further conveying devices, such as drift conveyors, which can also consist of belt conveyors. A scraper chain conveyor generally has at both its ends a chain drive, with which the endless scraper chain revolving in the upper strand and lower strand is driven with the scrapers which are connected thereto.\nA scraper for a so-called double-middle scraper chain conveyor, having a scraper web, a chain attachment region and scraper wings formed integrally onto the scraper web is known from DE 102 25 341 C1, to which is incorporated by reference herein in supplementation of the present disclosure.\nIn the prior art, ambitions to minimize the forces which shall be transmitted from the respective drives into the revolving scraper chain, in order to be able to reduce the capacities installed at the drives or, with the same capacity, transport higher quantities with the same scraper chain, have long existed. In EP 1 276 685 B1, it is proposed to this effect to make the scraper chain and/or the scrapers consist at least partially of a plastic which is self-extinguishing, flame-retardant and has anti-static properties, together with high mechanical strength. Through the use of plastic as the material for making the scraper chain and/or the scraper, a reduction in the weight of the scraper is intended to be achieved, whereby the drive capacity to be installed can be designed smaller or a larger conveying volume can be handled. According to one illustrative embodiment of EP 1 276 685 B1, the scraper is intended to consist of a metallic sheath, the interior of which is filled with a plastic and, according to another illustrative embodiment, the scraper is intended to be made fully of a suitable plastic, where possible the scraper chain also.\nFrom WO 2004/048234 A1 it is known to make the entire scraper of plastic, the plastics scraper being intended to be connected to the metal chain links of a metallic scraper chain. In order to achieve a close attachment between scraper and chain, the fully plastics scraper is intended to be cast directly onto the chain links of the scraper chain. A drawback with this is that the connection can only be established above ground, so that the entire scraper chain would have to be brought from below ground to above ground in order to exchange such scrapers. An exchange of individual scrapers in the installed state of the scraper chain in the scraper chain conveyor is impossible.\nDE 23 53 005 A1 discloses scrapers for a single-chain conveyor, in which a link of the simple chain is inserted between the top web and the top side of a chain clamp in order to connect the scraper to the simple chain. In order to reduce noise and wear, the sliding surface of the scraper consists of plastic. To this effect, the bottom side of the chain clamp can be coated with a plastics sliding surface, or it is provided with a plastics moulding which is anchored in a dovetailed transverse groove running transversely to the motional direction. In addition, the scraper wing ends can consist of exchangeable plastics caps, which are respectively fastened to a holding portion on the top web by means of a clamping sleeve. The plastics caps are provided on their bottom side with slide beads.\nFor double-middle scraper chain conveyors, it is known from WO 2004/065270 to have the top web and the chain clamp consist of different materials in order to minimize the weight of the scraper. As materials for the top part, plastic, aluminium or ceramics are proposed in order to obtain a top part which is as lightweight as possible. The scraper ends can also consist of end caps, which are movably fastened to the top part by means of clamping sleeves and a tongue and groove connection, for which holding portions are configured on the top web, which holding portions form the tongue and engage in a horizontal groove on the inner sides of the end caps."}
{"text": "1. Field of the Invention\nThe present invention relates to a club-head for a golf club and, more particularly, to an improvement of a club-head having a hitting portion which includes a main body made of a fiber-reinforced plastic and a metal sole member integrally fixed to the main body along the underside of the main body.\n2. Description of the Related Arts\nRecently, a club-head has been used in which a hitting portion provided at the heel thereof with a hosel or neck portion for connecting a shaft consists of a main body made of a fiber-reinforced plastic for hitting a ball at the front surface thereof and a sole member made of a metal is integrally joined to the main body along the underside thereof. The club-head having such a construction has a drawback in that the main body is apt to be exfoliated from the sole member upon impact with a golf ball.\nU.S. patent application No. 840,795 filed by the present applicant on Mar. 18, 1986, discloses one kind of the above-mentioned club-head wherein the back weight plate made of a heavy metal is attached to or embedded in the rear surface of the main body above the sole member. It has known that the club-head having such a construction can increase a flight distance of a golf ball hit by the head because the heavy back weight plate can most efficiently serve the main body to effect a repulsion action on the golf ball. But, the club-head having such a construction also has the drawback in that the main body is apt to be exfoliated from the sole member upon impact with a golf ball.\nU.S. patent application No. 854,626 filed by the present applicant on Apr. 22, 1986, discloses another kind of the above-mentioned club-head wherein the weight member is embedded in the main body and connected to the sole member through one or more connecting members. In this construction, when the heavy weight member is arranged between the front and rear surfaces of the hitting portion of the club-head, the thickness of the main body between the front surface of the hitting portion of the club-head and the weight member is decreased, and thus, the repulsion action of the head against the golf ball and depth of the center of gravity of the head are decreased. On the other hand, when the heavy weight member is located at the rearmost position of the hitting portion of the club-head, the head has a drawback in that the main body is apt to be exfoliated from the sole member."}
{"text": "Today's modern aircraft use hermetically sealed fire extinguishers that are opened, or activated, by direct explosive impingement energy. With reference to Prior Art FIG. 1, the device which provides the explosive energy is typically called a cartridge 120, or squib. The impingement energy is focused on a dome-shaped hermetic burst disc 110 such that the burst disc will rupture as a result of the impingement. The burst disc material used is typically fabricated from corrosion resistant steel.\nTypically, the cartridge 120 is retained in a discharge head 130 in such a manner that it directly faces the burst disc 110 assembly. The discharge head 130 is attached to the outlet of the fire extinguisher and is typically used to direct the flow of extinguishing agent to an aircraft interface, such as plumbing or tubing, which directs the agent to the desired location. A filter screen 150 is located within the discharge head to catch any large burst disc fragments created as a result of the explosive impingement energy."}
{"text": "Overall implementation of a keypad for a mobile device, i.e., the final package, including domes, backlight and keyset, should provide the required visual functionality as well as sufficient tactile feedback for the user when pressing an individual key. In view of thin keypads, the tactile feedback can be obtained using elastic light guide, dome pin lead throughs (holes above domes), or dome pins integrated with a plastic light guide (flexible but not elastic). Of these, sufficient optical quality can be obtained using either the first or the last approach.\nSpatially selective out-coupling of light from a planar light guide plate for keypad illumination is an important desired feature. It means that simultaneous control of out-coupling efficiency and spectra of the out-coupled light needs to be provided to enable modification of the visual appearance (intensity, colored/white) of the individual keys. In addition, a tactile feedback when pressing an individual key should not be compromised by the illumination system.\nVarious types of out-coupling means based on surface structuring of the planar light guide exist. Spatial modulation of the out-coupling can be obtained by limiting/modulating the spatial structuring to/within separate regions (key areas) on the light guide.’\nMoreover, colored appearance can be typically obtained using white light out-coupling and ‘painted’ keysets on a keymat. In cost sensitive applications, such as mobile phones, the utilized ‘paints’ typically have large variety of optical quality depending on the utilized manufacturer and can cause significant efficiency loss due to poor optical characteristics.\nVarious difficulties appear in keypad utilizing light guides for illumination and usually requiring several different manufacturing steps in keypad and light guide fabrication. Manufacturing of keypad's dome pins is usually realized by injection molding. The mechanical key function is realized by dome pins through holes in the illumination light guide. This can lead to energy losses in the illumination."}
{"text": "1. Field of the Invention\nThe present invention relates to apparatus, method and system for providing a dusting attachment device, which is installed on a suction-type vacuum cleaner and functions to equalize air flow current while providing a favorable dusting or dislodging surface.\n2. Background Information\nThose references found in the process of a U.S. patent search, having some relation in the prior art to technology attempting to address dusting and cleaning functions, or marshalling of debris in relation to another cleaning assembly, of those references found, include United States Patents to Miller, U.S. Pat. No. 5,123,142; Hamilton, U.S. Pat. No. 2,198,339; Palomino, Jr., U.S. Pat. No. 5,074,008; Smith, Jr., U.S. Pat. No. 4,972,541; Farmer, U.S. Pat. No. 4,198,727; Sorenson, U.S. Pat. No. 5,692,263; El Omary, U.S. Pat. No. 5,115,765; and Na, U.S. Pat. No. 5,652,997.\nMiller \"\"142 discloses a Vacuum Operated Cleaner Attachment. Disclosed therein is a cleaner attachment for use with a common vacuum cleaner having a group of duck feathers which are longitudinally arranged around half of its cylindrical tube. Miller\"\"s tube has the same cross-section throughout its entire length. The tube has a side wall and an end wall, and communicates, as claimed, with a vacuum source through its hollow internal passageway. This reference teaches a special securing means in the form of a wire overlying the quills of its feathers and heat shrinkable material overlying the wire holding the feathers. The feathers are each secured at one basic location, only, on the tube, and face the same longitudinal direction. A mustering of openings, having no discernible or disclosed pattern are provided through the side wall of Miller\"\"s tube member beneath the one-directional feathers. Specifically, as disclosed, illustrated or claimed; Miller\"\"s tube-side wall openings are beneath the longitudinal feathers and are said to be directing air flow laterally through the feathers, where the openings communicate with the passageway within the tube and serve to distribute generally equalized lateral air flow directly through the area of the feathers into the portion of the tube beyond the points of attachment of the feathers.\nMiller does not illustrate or claim a conical or pyramidal central core, a functional and discernible patterning of air channels, or the multiple and diverse connection sites of feathering members at interstitial spaces between air channels throughout the majority of perimeter or circumferential space available on a core member. Nor does Miller set forth any of the configurational arrays of feathering set forth by the present invention; only one longitudinal arrangement of straight duck feathers secured at one basic location only; and dimensioned at basically one quill diameter and length of feather.\nPalomino \"\"008 discloses a Dust Mop Attachment For Vacuum Cleaners. In its broadest teaching Palomino discloses and claims a mop attachment to a nozzle of a vacuum cleaner, having no feathering, and no array or patterning, or plurality, of air channels between or amongst which feathering or a reasonable equivalent is attached. Rather, only, a mop having a frame member, a mop element with one centrally located hole for receiving a vacuum cleaner nozzle, and means for detachable securing the mop element to its frame member and for detachably securing the cleaner nozzle to the mop element. elements of the mop simply surround one basic air hole, in vast difference with the present invention, where a diversity and plurality of feather members are attached to interstitial areas between a plurality and array of air channels deployed about a cone-like or pyramid-like core unit.\nThe Hamilton \"\"339 reference issued in 1940 discloses simply a xe2x80x9cbrush,xe2x80x9d especially adapted for cleaning foreign matter from radiators, Venetian blinds and other such devices with parts difficult to reach for cleaning. The xe2x80x9cbrushxe2x80x9d is disclosed as being adaptable to a vacuum cleaner, and shaped in its main body for working between xe2x80x98slatsxe2x80x99 of Venetian blinds or within and about the various tubes of a radiator, to loosen or dislodge, with the force and structure of a xe2x80x9cbrush,xe2x80x9d the dirt and foreign matter encountered, for vacuum cleaning removal. For this purpose, Hamilton\"\"s sole disclosure and claim address: xe2x80x9cA brush comprising an elongated hollow body having one end [its connecting end] reduced to form a neck for insertion into a suction hose and provided at the opposite end with an arcuate curvaturexe2x80x9d having an opening cut on a bias. Hamilton\"\"s hollow body is xe2x80x9cof substantially elliptical shape in cross section to provide opposed edges and opposed convex walls,xe2x80x9d in almost opposite distinction to the structure of the present invention. The xe2x80x98Brushxe2x80x99 of this reference is further provided with openings on the brush body which are disclosed as essentially being two lengthwise oriented lines of openings, one on each of the body\"\"s two opposed convex walls, as shown in cross-section in Hamilton\"\"s FIG. 2 drawing. Also, Hamilton discloses and claims xe2x80x9ctufts of bristles secured on the edges and opposed convex walls of the [brush] body;xe2x80x9d rather than a circumferentially oriented array and patterning of a number of air apertures on each arced or planar surface, with an interstitially attached number of feathering components about all such areas of a conical or pyramidal core unit. Hamilton does not teach a plurality of air channels in a diverse or patterned array about its available spacing; and does not teach the placement of feathering, or even its own xe2x80x98brush bristles,xe2x80x99 in areas between such air apertures or holes (interstitial areas) on either of the two opposing elliptical faces that Hamilton does teach.\nSmith \"\"541 discloses a feather duster having its own built-in fan assembly, dust catching assembly and turbulence-inducing gap. Disclosed is a xe2x80x9cmultiplicity of feathersxe2x80x9d each of which is attached at one end, only, of the device, pointing in one direction relative to the device; and without a positional or functional relationship to a patterning, array and plurality of air apertures about a core unit attachment, as in the case of the present invention.\nThe remaining references found appear to be of less significance than those discussed previously herein. They pertain to a baseboard duster for vacuum cleaners (\"\"727); a delicate dusting vacuum tool having outer holes sunk within grooves and no feathering components or reasonable structural similarity or equivalent (\"\"263); a flea blower vacuum for pets (\"\"765); and an auxiliary suction tool for use in a vacuum cleaner (\"\"997).\nNone of the references found in the prior art specifically illustrate or disclose the vacuum dusting attachment device of the present invention. Nor is the invention obvious in view of any of the prior art references listed. In addition, all of the relevant prior art heretofore known suffer from a number of disadvantages.\nNone of the prior art references adequately address the problem of marshaling dust and particle debris to a central body for suction and evacuation.\nAlso, the prior art has not provided an accurate and flexible feathering arrangement for efficiently dislodging dust and like particles, but relied upon hard or substantially inflexible bristles or one-direction, hard, feathering of one basic length and dimensioning, secured to point in only one direction, and secured at one location only or at limited locations on inferior designed core bodies.\nAdditionally, none of the prior art devices in this area of technology have employed any configuration, or dimensioning or sizing, of feathering, except those of conventional dusters utilized in the art without vacuum cleaning devices. Nor has the prior art addressed the utilization of conical or pyramidal core units, uniformly tapered from their base end to their apex end, providing a diverse patterning of air channels over about two-thirds (⅔\"\"s) of its length, and at the tip area, or providing a plurality of feather attachment sites at the spaces between each of a number of air channels on all faces of a core unit, for greater equalization of air currents and their effects more closely to a work site being dusted or contacted for dislodging particles therat.\nFurther, none of the prior art devices appear to have a central body or core which is constructed of malleable and flexible soft metal, alloy, polymer or composite; and which can be dimensionally adjusted by virtue of such construction materials, while retaining structural and positional integrity.\nThese and other disadvantages, structurally and functionally, of the prior art, will become apparent in reviewing the remainder of the present specification, claims and drawings.\nAccordingly, it is an object of the present invention to provide a substantially improved dusting attachment device for use in interaction with a vacuum cleaner, vacuum/suction cleaning means or system, or blowing cleaner, or vacuum, system; having improved air-equalizing and particle or dust dislodging features.\nIn this regard it is a further object of the present invention to provide an improved conical or pyramidally configured and shaped core unit, tapering from its base end to its apex end and having improved positioning, patterning and dimensioning of air apertures; and improved feathering means, of a number of novel and improved configurations and dimensioning, with individual feather components attached to to its novel core unit at locations between the air apertures throughout about two-thirds (⅔\"\"s) of the core\"\"s available overall spacing or length.\nIt is yet a further object of the present invention to provide a dusting attachment device having a core unit which is fabricated of soft, malleable and flexible metals, metal alloys, polymers, composite wood or metal-polymer materials, and positionable fabric material; each having positional integrity as adjusted or fitted; for ease in fitting, and size or dimensioning adjustment, in interaction with various cleaning system lines; and for simplicity and cost effectiveness in manufacturing and constructing the present invention. Yet a further related object, in this regard, is to provide cost effectiveness in the manufacturing process in utilizing many more diverse types and sizes of feathering and the convenient, efficient and inexpensive means of attaching the invention\"\"s feather components in spacing between its air apertures by means of utilizing a diverse number of conventionally available glues, cements, polymers, caulks or like fluids.\nIt will, therefore, be understood that substantial and distinguishable structural and functional advantages are realized in the present invention over the prior art devices; and that the present invention\"\"s simplicity and diversity of structure, diverse utility, and broad functional applications serve as important bases of novelty and distinction in this regard.\nThe foregoing and other objects of the invention can be achieved with the present invention, device, system and method which is a dusting attachment device for use in interaction with a cleaning system having a suction or blowing means and an line in communication therewith. The invention is provided with a core unit having apex and base ends and first, second and third portions therebetween, respectively, from the apex end to the base end, the core unit defining therewithin an air channel extending from the apex end, where it is enclosed thereat, to the base end, where it extends therethrough, and defining by virtue thereof an inboard surface adjacent to the air channel and an outboard surface spaced from the inboard surface. The apex end and the first and second portions of the core unit further define, in a patterning throughout a plurality of air apertures, extending through the outboard and inboard surfaces, and defining solid interstitial space portions between each of the air apertures on the apex end and on each of the outboard and inboard surfaces.\nThe invention is further provided with feathering means being fixedly attached to the solid interstitial space portions of the outboard surface of the apex end and the first and second portions of the core unit, the feathering means defining a configurational shape.\nThe dusting attachment can, therefore, be installed along the air channel of the base end, at the third portion of the core unit, in interaction with a line in communication with a suction or blowing means of a cleaning system, for functional use in equalizing the suction or blowing area between the apex end, the second portion and adjacent areas of the core unit, while dislodging dirt, dust, particles or debris in a work area."}
{"text": "1. Field of the Invention\nThe present invention relates to an information processing system, a mobile terminal, a server apparatus, a method for processing information, and a non-transitory computer readable storage medium.\n2. Description of the Related Art\nA technique has conventionally been known that collects motion information indicating motions of an indefinite number of people such as pedestrians, thereby detects the occurrence of an event (an abnormality, an intriguing event, or the like) to which a plurality of people pay attention, and estimates an area in which the event is occurring.\nFor example, a technique is known that collects positional information of mobile terminals carried by pedestrians or other people, thereby quickly detects a situation in which a plurality of people are stopping because of the occurrence of an abnormality, and identifies an area in which the abnormality is occurring.\nAnother technique is known that installs an imaging apparatus at a predetermined position, detects sight lines of a plurality of people using image information obtained by imaging an indefinite number of people such as pedestrians to identify an area on which the sight lines are focused, and thereby detects that an event has occurred in that area.\nHowever, simply by collecting the positional information of the mobile terminals carried by pedestrians or other people as described above, it is difficult to estimate the area in which the abnormality is occurring with high precision. This is because when people are looking at an abnormality at a distance, the positional information of positions at which the people are stopping does not necessarily represent the area in which the abnormality is occurring, for example.\nIn contrast, the method that identifies the area on which the sight lines are focused to detect the area in which the abnormality is occurring can avoid this problem. However, there is a problem in that when sight lines of an indefinite number of people are attempted to be monitored to constantly detect an area in which an abnormality is occurring, a huge amount of data is required to be processed at high speed, causing calculation costs."}
{"text": "Image processing systems are at the heart of digital image revolution. These systems process the captured digital image to enhance the clarity and details of the image using image processing algorithms. Such algorithms result in images that are substantially more accurate and detailed than previously achieved using older analog methods.\nThere remains, however, a substantial difference between how an image is perceived by a person and an image captured and reproduced on a display medium. Despite the improvements gained by conventional digital image processing systems, such systems are still deficient in reproducing an image with the same level of detail, color constancy, and lightness of an actual scene as the eye, brain, and nervous system of a human being. This is due in part because the human nervous system has a greater dynamic range compression than is available on current digital systems. Dynamic range compression refers to the ability to distinguish varying levels of light.\nThe human eye has a dynamic range compression of approximately 1000:1, which means that the human eye can distinguish approximately 1000 levels of light variations. By way of contrast, digital image systems typically use only eight bits per pixel which allows for a dynamic range compression of only 255:1. As a result, a digital image reproduced as a photograph would have far less detail in the darker and brighter regions of the photograph as compared to the actual scene perceived by a viewer.\nMany techniques have been developed to compensate for this lighting deficiency. These techniques can be separated into two broad categories: (1) power law or nonlinear techniques (“non-linear transforms”); and (2) retinex techniques. Each have their respective limitations however.\nNon-linear techniques use a non-linear relationship to expand one portion of the dynamic range while compressing another. These techniques generally enhance details in the darker regions at the expense of detail in the brighter regions. One problem associated with these non-linear systems is that they provide greater distinctions between pixels regardless of what lightness value a pixel may have been originally assigned. This results in the brighter areas which already have a “washedout” appearance to become even more washed-out. Although these techniques result in better details in the darker regions, they do so at the expense of the brighter areas of the digital image. Further, these methods cannot handle contours and abrupt boundaries well.\nRetinex techniques variations increase or decrease the luminance value for a pixel based on the luminance values of surrounding pixels. These techniques are particularly useful for enhancing boundaries between lighter and darker regions of an image. However, such techniques are unsatisfactory for a number of reasons. In one technique, a large uniform zones in the image are grayed out (i.e., shading effect). In another technique a shift in color occurs in some images (i.e., color distortion) and is computationally intensive."}
{"text": "The present invention relates to a microfabrication technique which may be advantageously used, for example, for the manufacture of magnetic bubble memory devices, semiconductor devices, and the like. More particularly, the present invention relates to a method for forming a pattern of metal elements arranged with very small gaps therebetween on a substrate of a device.\nIn well-known field access types of magnetic bubble memory devices, a permalloy pattern for the propagation of magnetic bubbles (permalloy propagation pattern) is formed on the bubble-supporting magnetic layer of a substrate. The permalloy propagation pattern is constituted by a plurality of pattern elements made of thin films of permalloy. These permalloy pattern elements are chevron-shaped or half disc-shaped, for example, and are arranged with very small gaps therebetween. In a magnetic bubble memory device using 2 .mu.m diameter bubbles, the size of the gap is approximately 1 .mu.m.\nThe pattern of metal elements such as the permalloy propagation pattern described above is conventionally formed by using photolithographic techniques. That is, first, a pattern forming layer of metal is formed on a substrate of a device. Next, the pattern forming layer is coated with a layer of photosensitive resist material (photoresist). An optical exposure process is used to form a resist pattern corresponding to the pattern of metal elements to be formed from the layer of photoresist. Thereafter, an etching process is used to form the pattern of metal elements from the pattern forming layer.\nHowever, the smallest size of the resist pattern which can be formed on the pattern forming layer by using the photolithographic technique is approximately 1 .mu.m. Therefore, according to the conventional method, it is impossible to form a pattern of metal elements with gaps of a size smaller than 1 .mu.m.\nIn recent years, there have been a number of proposals for realization of a high density magnetic bubble memory device using 1 .mu.m diameter bubbles. Such a device requires a permalloy propagation pattern with gaps of a size less than 0.5 .mu.m. In order to satisfy this requirement, the electron-beam lithographic technique or X-ray lithographic technique is used. However, electron-beam systems and X-ray systems are large in size and high priced, thereby increasing the overall cost of the equipment."}
{"text": "1. Field of Endeavor\nThe present invention relates to measuring and more particularly to optically measuring interior cavities.\n2. State of Technology\nU.S. Pat. No. 5,432,543 issued Jul. 11, 1995 to Jun Hasegawa et al for endoscopic image processing device for estimating three-dimensional shape of object based on detection of same point on a plurality of different images provides the following state of technology information: “When the relative arrangement of the imaging means disposed with overlap in a visual field is known, various methods are known which estimate the shape of an object from the images formed, i.e., for estimating the shape from stereo-images. In recent years, some methods have also been proposed for determining a three-dimensional structure from information about the movement of an object. The methods employ the technique of estimating the relative movement of imaging means from a plurality of images.”\nU.S. Pat. No. 6,009,189 issued Dec. 28, 1999 to David F. Schaack for apparatus and method for making accurate three-dimensional size measurements of inaccessible objects provides the following state of technology information: “In the past several decades, the use of optical endoscopes has become common for the visual inspection of inaccessible objects, such as the internal organs of the human body or the internal parts of machinery. These visual inspections are performed in order to assess the need for surgery or equipment tear down and repair; thus the results of the inspections are accorded a great deal of importance.”\nU.S. Pat. No. 6,459,481 issued Oct. 1, 2002 to David F. Schaack for simple system for endoscopic non-contact three-dimensional measurement provides the following state of technology information: “Making accurate dimensional measurements of objects viewed through endoscopes is important to aerospace as well as other industries in which expensive equipment must undergo periodic internal inspections to maintain safe operation. Such measurements also have medical applications, where the internal condition of a patient is evaluated prior to or during surgery by viewing through an endoscope. The fundamental problems in making an accurate measurement through an endoscope are that the magnification of the image varies rapidly with the range of the object, and that objects of interest (defects) lie on surfaces which are curved in three dimensions; thus the magnification varies from one point on the object to another. What is needed is a fully three-dimensional measurement, that is, one which determines the depth, as well as the height and width, of an object. Endoscopes are long and narrow optical systems, typically circular in cross-section, which can be inserted through a small opening in an enclosure to give a view of the interior. They almost always include a source of illumination which is conducted along the interior of the scope from the outside (proximal) end to the inside (distal) end, so that the interior of the chamber can be viewed even if it contains no illumination. Endoscopes are divided into two basic types: these are the rigid “borescopes” and the flexible “fiberscopes” or “videoscopes.”"}
{"text": "1. Technical Field\nThe embodiments described herein are related to the field of golf balls and, more particularly, to a spherically symmetrical golf ball having a dimple pattern that generates low-lift in order to control dispersion of the golf ball during flight.\n2. Related Art\nThe flight path of a golf ball is determined by many factors. Several of the factors can be controlled to some extent by the golfer, such as the ball's velocity, launch angle, spin rate, and spin axis. Other factors are controlled by the design of the ball, including the ball's weight, size, materials of construction, and aerodynamic properties.\nThe aerodynamic force acting on a golf ball during flight can be broken down into three separate force vectors: Lift, Drag, and Gravity. The lift force vector acts in the direction determined by the cross product of the spin vector and the velocity vector. The drag force vector acts in the direction opposite of the velocity vector. More specifically, the aerodynamic properties of a golf ball are characterized by its lift and drag coefficients as a function of the Reynolds Number (Re) and the Dimensionless Spin Parameter (DSP). The Reynolds Number is a dimensionless quantity that quantifies the ratio of the inertial to viscous forces acting on the golf ball as it flies through the air. The Dimensionless Spin Parameter is the ratio of the golf ball's rotational surface speed to its speed through the air.\nSince the 1990's, in order to achieve greater distances, a lot of golf ball development has been directed toward developing golf balls that exhibit improved distance through lower drag under conditions that would apply to, e.g., a driver shot immediately after club impact as well as relatively high lift under conditions that would apply to the latter portion of, e.g., a driver shot as the ball is descending towards the ground. A lot of this development was enabled by new measurement devices that could more accurately and efficiently measure golf ball spin, launch angle, and velocity immediately after club impact.\nToday the lift and drag coefficients of a golf ball can be measured using several different methods including an Indoor Test Range such as the one at the USGA Test Center in Far Hills, N.J., or an outdoor system such as the Trackman Net System made by Interactive Sports Group in Denmark. The testing, measurements, and reporting of lift and drag coefficients for conventional golf balls has generally focused on the golf ball spin and velocity conditions for a well hit straight driver shot—approximately 3,000 rpm or less and an initial ball velocity that results from a driver club head velocity of approximately 80-100 mph.\nFor right-handed golfers, particularly higher handicap golfers, a major problem is the tendency to “slice” the ball. The unintended slice shot penalizes the golfer in two ways: 1) it causes the ball to deviate to the right of the intended flight path and 2) it can reduce the overall shot distance.\nA sliced golf ball moves to the right because the ball's spin axis is tilted to the right. The lift force by definition is orthogonal to the spin axis and thus for a sliced golf ball the lift force is pointed to the right.\nThe spin-axis of a golf ball is the axis about which the ball spins and is usually orthogonal to the direction that the golf ball takes in flight. If a golf ball's spin axis is 0 degrees, i.e., a horizontal spin axis causing pure backspin, the ball will not hook or slice and a higher lift force combined with a 0-degree spin axis will only make the ball fly higher. However, when a ball is hit in such a way as to impart a spin axis that is more than 0 degrees, it hooks, and it slices with a spin axis that is less than 0 degrees. It is the tilt of the spin axis that directs the lift force in the left or right direction, causing the ball to hook or slice. The distance the ball unintentionally flies to the right or left is called Carry Dispersion. A lower flying golf ball, i.e., having a lower lift, is a strong indicator of a ball that will have lower Carry Dispersion.\nThe amount of lift force directed in the hook or slice direction is equal to: Lift Force*Sine (spin axis angle). The amount of lift force directed towards achieving height is: Lift Force*Cosine (spin axis angle).\nA common cause of a sliced shot is the striking of the ball with an open clubface. In this case, the opening of the clubface also increases the effective loft of the club and thus increases the total spin of the ball. With all other factors held constant, a higher ball spin rate will in general produce a higher lift force and this is why a slice shot will often have a higher trajectory than a straight or hook shot.\nTable 1 shows the total ball spin rates generated by a golfer with club head speeds ranging from approximately 85-105 mph using a 10.5 degree driver and hitting a variety of prototype golf balls and commercially available golf balls that are considered to be low and normal spin golf balls:\nTABLE 1Spin Axis, degreeTypical Total Spin, rpmType Shot−302,500-5,000Strong Slice−151,700-5,000Slice01,400-2,800Straight+151,200-2,500Hook+301,000-1,800Strong Hook\nIf the club path at the point of impact is “outside-in” and the clubface is square to the target, a slice shot will still result, but the total spin rate will be generally lower than a slice shot hit with the open clubface. In general, the total ball spin will increase as the club head velocity increases.\nIn order to overcome the drawbacks of a slice, some golf ball manufacturers have modified how they construct a golf ball, mostly in ways that tend to lower the ball's spin rate. Some of these modifications include: 1) using a hard cover material on a two-piece golf ball, 2) constructing multi-piece balls with hard boundary layers and relatively soft thin covers in order to lower driver spin rate and preserve high spin rates on short irons, 3) moving more weight towards the outer layers of the golf ball thereby increasing the moment of inertia of the golf ball, and 4) using a cover that is constructed or treated in such a ways so as to have a more slippery surface.\nOthers have tried to overcome the drawbacks of a slice shot by creating golf balls where the weight is distributed inside the ball in such a way as to create a preferred axis of rotation.\nStill others have resorted to creating asymmetric dimple patterns in order to affect the flight of the golf ball and reduce the drawbacks of a slice shot. One such example was the Polara™ golf ball with its dimple pattern that was designed with different type dimples in the polar and equatorial regions of the ball.\nIn reaction to the introduction of the Polara golf ball, which was intentionally manufactured with an asymmetric dimple pattern, the USGA created the “Symmetry Rule”. As a result, all golf balls not conforming to the USGA Symmetry Rule are judged to be non-conforming to the USGA Rules of Golf and are thus not allowed to be used in USGA sanctioned golf competitions.\nThese golf balls with asymmetric dimples patterns or with manipulated weight distributions may be effective in reducing dispersion caused by a slice shot, but they also have their limitations, most notably the fact that they do not conform with the USGA Rules of Golf and that these balls must be oriented a certain way prior to club impact in order to display their maximum effectiveness.\nThe method of using a hard cover material or hard boundary layer material or slippery cover will reduce to a small extent the dispersion caused by a slice shot, but often does so at the expense of other desirable properties such as the ball spin rate off of short irons or the higher cost required to produce a multi-piece ball."}
{"text": "There is increasing interest in IP telephony to help lower costs and enable new services. Many enterprises and call centers are adopting IP telephony over their converged IP infrastructure and many multi-site corporations are using Voice over IP (VoIP) for their intra- and inter-site communication.\nWith the use of VoIP for mission-critical business applications, it is important to evaluate and improve the reliability and quality of VoIP calls. Ideally, a VoIP call should be as reliable as a traditional circuit-switched phone call. However, when monitoring and evaluating the quality of a VoIP call (especially over wide area links), it is necessary to deal with the inherent packet losses, delays, and jitter associated with IP networks, which are not encountered in traditional circuit-switched networks. Even though IP networks are largely self-healing for network faults, and many enterprise networks are engineered to have redundant links or paths between sites, today's IP networks are not engineered to react to performance degradations at the timescales needed for voice. For example, recent studies show that while there is acceptable performance within some service provider networks, many backbone paths still have poor VoIP performance and network faults cause problems.\nTo provide a robust VoIP infrastructure, it is important to rapidly detect performance degradations and faults. This detection is complicated by several factors. On a per-connection basis, for example, there are natural silence periods in VoIP calls during which packets are not transmitted by a source (e.g., when a participant in a call is listening rather than speaking). Consequently, while monitoring a VoIP call (e.g., on the receiving side) it is necessary to distinguish between gaps that occur due to natural speech silences and perhaps speech compression, and the gaps that occur due to packet loss, delay, and jitter in the IP network. Furthermore, although detecting problems can certainly help to alert a network manager, it would be particularly useful if the network could react to a detected problem and route around it.\nAccordingly, there is a need for techniques for performing rapid fault detection and recovery in communication networks such as IP telephony networks, particularly those that provide VoIP applications."}
{"text": "The present invention relates to temperature regulation of the human body, in particular, to the use of heat pipes in garments, blankets and pads to warm or cool the human body for activities in extreme temperature environments and for treatment of medical conditions.\nThe need to regulate body temperature exists where human activities are conducted in extreme temperature environments. Very cold environments are encountered by undersea divers, inhabitants of polar regions, pilots, mountain climbers and snow skiers, as well as by individuals who work in more pedestrian cold environments such as refrigerated containers. Very hot environments are encountered by those in diverse occupations such fire fighters, foundry workers and armored tank drivers.\nThe need for body temperature regulation also exists in the treatment of medical conditions. Individuals with impaired body temperature regulation systems, such as the elderly, stroke patients, and patients with spinal cord injuries or multiple sclerosis, may be subject to chronic hypothermia or hyperthermia. Victims of exposure, as well as patients having fever, dehydration, infection, and drug reactions may suffer temporarily from hypothermia or hyperthermia, and may require treatment including body temperature regulation.\nDeliberate inducement of whole body hyperthermia has been found to be useful as a treatment for cancer, and has met with success. However, side effects range from fatigue to occasional mortality. Injury to temperature-sensitive body parts, such as the liver, has also been associated with whole body hypothermia. Because of these side effects, and the high labor costs associated with monitoring such treatment, regional, rather than whole body, hypothermia has found favor in the treatment of cancer. Inducement of regional or localized hyperthermia has, thus, also been used in the treatment of cancer, as well as in the treatment of joint diseases, such as capsulitis, tendinitis and lower back pain.\nVarious means have been developed to warm and cool the human body to achieve body temperature regulation needed for such activities and medical treatments. Cold suits, such as the garment shown by Jennings, U.S. Pat. No. 3,289,748, include circulating fluid systems powered by a pump or other means for circulating. See also, Macdonald et al, U.S. Pat. No. 4,807,447. In the garment of Jennings, serpentine tubes extend along the arms and legs of the wearer and converge at inlet and outlet manifolds. Cold suits, however, have the disadvantage of being heavy and are typically restrictive, causing exhaustion and cardiopulmonary distress in the wearer.\nHyperthermia blankets are also known wherein a single phase fluid loop winds back and forth through the blanket from inlet to outlet to cool the user. A chiller attached to the outlet of the fluid loop discharges heat from the hyperthermia blanket. Such blanket systems, however, have problems maintaining a constant blanket temperature throughout, and are heavy due to the weight of the circulating fluid. A blanket including heat pipes provides more uniform heating, as shown by Sgroi, U.S. Pat. No. 4,094,357. In the blanket of Sgroi, the heat pipes converge at one side of the blanket to be cooled by a refrigeration device or, alternatively, heated by a thermal energy source.\nSuits and blankets having circulating fluid loops in similar configurations may also be used to provide heat to persons suffering from hypothermia. See, Macdonald et al and Sgroi. However, the same general disadvantages of weight, restriction of movement, and resulting exhaustion are present. Electric blankets are also well-known but often provide uneven heating, and subject the user to low-level electromagnetic radiation. Nonetheless, whether used for heating or cooling, such suits and blankets generally disadvantageously require external connections and/or require external power sources.\nNumerous devices have been developed for regional, therapeutic heat transfer. Jones, U.S. Pat. No. 4,706,672 discloses a passive thermal transfer device in which metal fibers woven into a mesh act as a heat sink to draw heat from an afflicted area for dissipation. Golden, U.S. Pat. No. 4,962,761 discloses a thermal bandage to heat or cool an area of the skin, wherein a circulating fluid is used to carry heat to or from the afflicted part. However, the need continues to exists for improved means for applying heat to specific areas of the human body for the treatment of disease which will avoid damage to adjacent tissues and structures."}
{"text": "While iron ore raw materials are used as the raw material for producing reduced iron as might be expected, recent price surge of iron ore raw materials due to strong demand has brought attention to recycling of iron raw material-containing dust produced in iron mills. In pig iron- and steel-making processes, powdery iron and iron oxide-containing dust blown up in blast furnaces, converters, melting furnaces, electric furnaces, and the like are recovered with dust collectors. The recovered powders (referred to as “steel-making dust” in Claims and the specification) are used as iron raw materials since they contain iron and iron oxides.\nIn a rotary hearth furnace, which is one of facilities for producing reduced iron, raw materials containing metal oxides and carbonaceous substances must be fed in the form of agglomerates to achieve homogeneous reduction reactions, and the agglomerates need to have a particular strength. However, when carbonaceous substances are incorporated, since carbon burns at a relatively low temperature, a technique of increasing the strength of the agglomerates by sintering, such as in the cases of sintered pellets and sintered ore, cannot be employed. Thus, the strength of the agglomerates has been enhanced by use of expensive binding agents (hereinafter also referred to as “binders”), such as starch and molasses.\nTypically, either one of a pellet-making facility or a briquette-making facility is used as the facility to produce agglomerates, and, in either cases, agglomerates having high strength and as uniform size as possible must be produced. However, it has been extremely difficult to produce agglomerates from fine raw materials mainly for two reasons described below. The first reason is that the fine powder dust raw materials have a small bulk density and many gaps and thus do not readily form high-strength agglomerates. If the strength of the agglomerates is low, the agglomerates may collapse under pressure and thus cannot be stored in large quantities, and cracking and the like may occur during transportation. The second reason is the difficulty of homogeneously mixing a binder, which is mixed into the fine raw materials to increase the strength of the agglomerates. This is because a binder having an adhesive function has viscosity, and this viscosity obstructs homogeneous mixing. In addition, the binder materials are not only expensive but also can cause clogging of materials during manufacture or feeding (e.g., in an intermediate hopper or the like) of the agglomerates since the binder materials have the adhesive function. Once clogging of the raw materials occurs, the supply of raw materials to the rotary hearth furnace or the like stops or a large quantity of raw materials are suddenly fed, which renders the operation of the reducing furnace unstable. Thus, the amount of the binder used is preferably as small as possible.\nIn the description below, methods for producing pellets and briquettes of fine dust raw materials that have been known in the art are explained by citing examples.\nPatent Document 1 teaches a method for producing pellets including kneading raw materials, which has been discharged from a raw material storage, by using a kneader and pelletizing the raw material, which includes grains containing a metal oxide and carbon, by using a pan pelletizer, as shown in FIG. 1 of Patent Document 1. In particular, a plurality of raw materials are fed onto a raw material conveyer from a plurality of material reservoir bins at a determined mixture ratio. The grain diameter of the raw material, the chemical composition, and the ratio of water to be mixed are determined. In particular, in order to adequately carry out reduction reactions, the ratio of the metal oxide to carbon is adjusted. The claims of Patent Document 1 disclose a method for producing pellets suitable for reducing furnaces, the method being characterized in that, in producing spherical pellets with a pan pelletizer from a powder containing a metal oxide containing steel-making dust and a carbon-containing powder, the powder contains 20 to 80% of grains having a diameter of 10 μm or less.\nPatent Document 2 teaches a method for producing a solidified product from steel-making dust as shown in FIG. 1 of this patent document. As shown in this drawing, this system includes a granulating step of mixing the dust mainly composed of iron and its oxides produced during steel making with a powder mainly composed of carbon and pelletizing the resulting mixture, a water impregnation step of impregnating the mixture granules with water, a solidifying step of placing the water-impregnated pellets into forming dies and pressure-compacting the pellets to form briquettes, a transfer step of transferring the briquettes as the raw material for a melting furnace, and a dezincing step of concentrating zinc in the dust and dezincing the dust in the concentrated state after obtaining the dust from the melting furnace and before the solidifying step. According to this method, since impregnation with water is conducted immediately before molding, the surfaces of the mixture granules are softened to promote deformation of granules during briquetting, and the adhesion strength between the granules can be increased as a result."}
{"text": "1. Field of the Invention\nThe present invention relates to a static random access memory having high-resistance elements made of polycrystalline silicon and used as load elements within memory cells, and more particularly, to a static random access memory provided with a means for detecting defective memory cells having leakage current, if any.\n2. Description of the Related Art\nAn E/R type static random-access memory cell (E/R type SRAM cell) is known. As is understood from the equivalent circuit diagram shown in FIG. 1, this memory cell comprises a flip-flop circuit. The flip-flop circuit is comprised of two inverters, each having one high-resistance element which is used as a load element, and one enhancement type MOS transistor which functions as a drive element. The input and output of the first inverter are connected to the output and input of the second transistor, respectively. More specifically, high-resistance elements R1 and R2 are coupled, at one end, to power-supply voltage V.sub.CC, and at the other end, to the drains of MOS transistors Q1 and Q2, respectively. The sources of MOS transistors Q1 and Q2 are ground voltage V.sub.SS. The gate of transistor Q1 is coupled to the drain of transistor Q2, and the gate of transistor Q2 is coupled to the drain of transistor Q1. Hence, high-resistance elements R1 and MOS transistor Q1 form an inverter, and high-resistance elements R2 and MOS transistor Q2 form another inverter. These inverters have their input and output connected to each other's output and input, respectively, whereby the inverters constitute flip-flop circuit F.\nAs is illustrated in FIG. 1, MOS transistors Q3 and Q4 are used as transfer gates for controlling data-writing and data-reading. MOS transistor Q3 is connected between memory node N1 of flip-flop circuit F and bit line BL. MOS transistor Q4 is coupled between memory node N2 of flip-flop circuit F and bit line BL. The gate of both MOS transistors Q3 and Q4 are coupled to word line WL. All transistors Q1 to Q4 are of enhancement type.\nThe flip-flop circuit F, which is the main component of the memory cell, is a bistable one. Therefore, two data items complemental to each other are stored in memory nodes N1 and N2. Assuming that data items, i.e., \"1\" and \"0\", are stored in nodes N1 and N2, respectively, transistor Q1 is off, and transistor Q2 is on. Since high-resistance elements R1 and R2 are linear passive elements, a current flows through high-resistance element R2 coupled to transistor Q2 which is on. The resistance of high-resistance element R2 is far higher than that of on-resistance of transistor Q2. Hence, the resistance of element R2 determines the value of the current flowing through transistors Q2. A similar current flows in each of the E/R type SRAM cells forming a static random access memory (SRAM), and the sum of the currents flowing through all E/R type SRAM cells determines the SRAM in its static condition.\nAlso well known are completely CMOS SRAM cells, each comprising six transistors, two of which are P-channel MOS transistors functioning as load elements. The two P-channel MOS transistors, which are used as load elements, are active elements. Therefore, the memory cell of FIG. 1 can be made of four MOS transistors and two resistors, and occupies but less area than a fully CMOS memory cell. However, this memory cell is disadvantageous in that a current flows via both high-resistance elements R1 and R2 as long as the cell in the static condition. To reduce this current flowing through elements R1 and R2, these high-resistance elements can be replaced with resistors which have higher resistances. However, when elements R1 and R2 are replaced by such resistors, the flip-flop F may operate unstably, inevitably destroying the data stored in the E/R type SRAM cell, as will be explained below.\nFIG. 2 is a table showing the various resistances which elements R1 and R2 of each E/R type SRAM cell of an E/R type SRAM must have to maintain the current, which flows through the cell in the static condition, at the typical value of 1 .mu.A, in accordance with the memory capacity of the E/R type SRAM. As may be understood from FIG. 1, memory nodes N1 and N2 exist in the PN junction between a P-type substrate and an N+ diffusion region which is the common drain of transistors Q1 and Q2. The reverse leakage current flowing in this PN junction is about 10.sup.-14 A. This value is equivalent to resistance of about 100 tera .OMEGA.. In the case of a 256 K-bit memory, the resistance of either high-resistance element, R1 or R2 is 100 times the resistance equivalent to the reverse leakage current flowing in the PN junction. In the case of a 1 M-bit memory, the resistance is 25 times the resistance which is equivalent to the reverse leakage current. The greater the ratio of the resistance of either high-resistance element, the larger the operation margin of each memory cell. As is evident from FIG. 2, this operation margin is inversely proportional to the capacity of the memory.\nWhat has been explained in the preceding paragraph only applies to the case where no abnormal leakage currents flow in the PN junction connected to both memory nodes N1 and N2. In actuality, both a defective leakage current and a contaminating leakage current are generated in the PN junction. When these abnormal leakage currents are far greater than the currents flowing via high-resistance elements R1 and R2, neither memory node can no longer hold the data item \"1\", and the flip-flop circuit F ceases to function, and hence the E/R type SRAM cell having this flip-flop circuit F cannot correctly function. Any SRAM chip, if containing such a defective memory cell, can no longer function at all. Needless to say, the presence of such defective memory cells reduces the yield of SRAM chips.\nOn the hand, when the abnormal leakage currents are nearly equal to the currents flowing through high-resistance elements R1 and R2, the voltage applied on the memory node storing the data item \"1\" falls to the value determined by the division of resistance achieved by high-resistance elements R1 and R2 and also by the resistance equivalent to the abnormal leakage current. Consequently, flip-flop circuit F becomes very unstable, and the memory cell, as a whole, inevitably operates unstably. Any SRAM which includes such a defective memory cell has neither a sufficient margin for power-supply voltage nor a sufficient margin for temperature changes. Unless the defective memory cell is detected, and any appropriate measures are taken for the cell, the data stored in this defective cell will be destroyed eventually.\nHere arises a problem. It is extremely difficult to detect a defective memory in the SRAM. As is known in the art, it takes a long period of time to complete the test for screening an inadequate operation margin or a defect which needs a long time to detect. In some cases, it is necessary to set the ambient temperature at various values to carry out this test successfully. In the worst case, the insufficient operation margin or the defect may not be found even if the test has been carried out for a considerably long time.\nFIG. 3 is an equivalent circuit diagram showing the flip-flop circuit of an E/R type SRAM cell which is identical to the SRAM cell shown in FIG. 1, except that memory nodes N1 and N2 have a leakage-current path each. Resistors Rj and Rj' shown in this figures are equivalent to the leakage-current paths. More precisely, resistance Rj exists only when a reverse leakage current of a normal value flows through the PN junction coupled to memory node N1, whereas resistance Rj' exists only when a leakage current of an abnormal value flows through the PN junction connected to memory node N2.\nThe typical temperature characteristics of resistors R1, R2, Rj, and Rj' are illustrated in FIG. 4. As can be understood from FIG. 4, the energy for activating high-resistance elements R1 and R2, each formed of a polycrystalline silicon layer, is great, and the temperature dependency of either high-resistance element is prominent. The resistance of resistor Rj is determined by the junction area of memory node N1, and is about 100 tera .OMEGA.. The temperature dependency of resistor Rj is less than that of high-resistance elements R1 and R2. Resistor Rj' has a resistance which is the sum of the resistance of resistor Rj and the resistance equivalent to the abnormal leakage current flowing through the PN junction. As is evident from FIG. 4, the line representing the temperature characteristic of resistor Rj' crosses the line representing the temperature characteristic of resistor Rj at a specific temperature T1 if the resistance of resistor Rj' is about hundreds of times lower than that of resistor Rj, said specific temperature T1 falling within the range in which the memory cell can correctly perform its function. In other words, resistor Rj' has the same resistance as elements R1 and R2 at temperature T1. Its has a resistance lower than that of high-resistance elements R1 and R2 at any temperature below T1, and a resistance higher than that of elements R1 and R2 at any temperature above T1.\nThe voltage at which memory nodes N1 and N2 hold data when the memory cell is in its static condition will now be calculated. The voltage V1(1) which node N1 requires to hold data \"1\" can be given: ##EQU1##\nThe voltage V2(1) which node N2 requires to hold data \"1\" can be given: ##EQU2##\nThe voltage V1(0) which node N1 requires to hold data \"0\" can be given: ##EQU3##\nThe voltage V2(0) which node N2 requires to hold data \"0\" can be given: ##EQU4##\nIn the above equations (1) to (4), Rn(off) is the resistance equivalent to the channel leakage current which flows when transistors Q1 and Q2 of flip-flop circuit F are off, Rn(on) is the resistance equivalent to the channel current which flows when transistors Q1 and Q2 are on. Any transistor has Rn(off) as high as 10.sup.14 .OMEGA. or more if operating correctly. Hence, equations (1) to (4) can reduce to: ##EQU5##\nIt is desirable that the voltages which each memory cell requires to store data \"1\" and data \"0\" be V.sub.CC and 0 V, respectively. Voltage V2(1), at which memory node N2 holds data \"1\", is determined by Rj'/(R2+Rj'), and is lower than V.sub.CC. The lower Rj' is than R2, the more does voltage V2(1) fall. When voltage V2(1) falls below the threshold voltage Vthn of each N-channel MOS transistor, Rn(on), one of the terms of equation (3), changes to Rn(off). Consequently, voltage V1(0), at which memory node N1 holds data \"0\", will rise to V.sub.CC with time constant of R1.CA (where CA is the capacitance present in memory node N1). That is, data \"0\" held by node N1 has changed to \"1\", whereas data \"1\" held by node N2 has changed to \"0\". This is the destruction of data. In terms of the temperature characteristics shown in FIG. 4, the data is destroyed at temperature T1 or any lower temperature. This destruction of data is, hence, called \"leakage induced cell destruction at low-temperature.\"\nA SRAM memory cell which has undergone low-temperature pause destruction can hardly be detected unless the SRAM chip is tested not only at room temperature, but also at other various temperatures. It takes an extremely long time to test the SRAM chip at various temperatures, for a memory cell, if any, which has undergone leakage induced cell destruction at low-temperature. Obviously, it is uneconomical to spend much time on testing each SRAM chip. Further, a SRAM cell repeats unstable operation several times until the data stored in it is destroyed completely. This unstable operation of the cell is very hard to detect within a short time of period. Accordingly, demands have been made for means which can text a wafer having a number of SRAM chips, to determine quickly whether or not each SRAM chip contains any defective memory cell.\nHitherto, in order to detect defective memory cells, if any, contained in a SRAM chip, the operation characteristics of all cells of the SRAM chip are externally examined. This method can hardly be successful and requires a very long time to test the SRAM chip."}
{"text": "1. Field of the Invention\nThe present invention relates to demodulation circuits for use in direct sequence modulation systems and, more specifically, to demodulation circuits for use in biphase balanced modulation systems.\n2. Description of the Prior Art\nDirect sequence modulation systems are known in the art. In such systems, an RF carrier signal is modulated by a digital code sequence whose bit rate is much higher than the information signal bandwidth. One of the preferred modulation formats for such systems is the biphase balanced modulation scheme in which the carrier is transmitted with one phase when the code sequence is a binary ONE and with a 180.degree. phase shift when the code is a binary ZERO.\nVarious techniques for detecting biphase modulated signals are known. In general, such techniques are coherent processes requiring that the receiver include a local code reference that is an accurate estimate of the received code modulation. Prior art techniques for providing this local code reference are effective but require relatively complex circuitry.\nOne of the better known demodulators for this type of system, for instance, is the Costas loop, as described in the book: \"Spread Spectrum Systems\" by R. C. Dixon and published by John Wiley & Sons, New York, in 1976. However, Costas loops are relatively complicated in that they require a pair of demodulators in phase quadrature with each other. When the loop is phase locked, one of the demodulators will be in quadrature with the incoming signal and its output will be basically zero. The reference associated with the other demodulator will demodulate the signals at the maximum value, which may be either positive or negative, depending upon the instantaneous phase of the modulation. Typically, the incoming signal is integrated for a period of one bit, after which time the output of the demodulator which is in phase with the signal is sensed to determine if it is positive or negative, whereupon either the direct or inverted error signal may be used. Thus the Costas loop typically requires integrated dump circuitry, decision making circuitry, and means for inverting the error signal.\nIn contrast, the circuit of the present invention typically requires only a phase-lock loop, a flip-flop, and an exclusive OR gate."}
{"text": "In conventional photolithography a photoresist is exposed to light through a mask. The photoresist is modified by the exposure in such a way that either the exposed or unexposed portions of the resist can be removed during subsequent development. Any photolithographic process has limitations, whereby there is a critical dimension below which features are too fine to be resolved. That resolution limit is a critical barrier in reducing the scale of integrated circuit devices.\nSelf-aligned double patterning is a technique for forming features having a finer pitch than would be possible by the direct application of a photolithographic process. Self-aligned double patterning involves forming a mandrel having line-shaped features. A spacer formation process is then used to form spacers on the sides of the mandrel features. The mandrel is then stripped leaving the spacers defining two sets of lines. A first set of lines corresponds to the mandrel. A second set of lines is formed between each adjacent pair of the mandrel's line-shaped features, the second set of lines being form between adjacent spacers."}
{"text": "Advances in electronics, batteries and packaging technology have led to continued increases in the number of mobile computing devices in service. These mobile computing devices often have an associated docking station that permits ready access to printers, modems, networks, and peripherals that are more humanly comfortable, amongst other connections and attributes. Though beneficial for their intended purpose, these docking stations are to some extent disadvantageous as discussed below.\nMobile computing devices tend to be small to facilitate efficient mobility. Docking stations tend to be configured for desktop use in which a mobile computing device is docked at the station for data transmission or for running the desktop environment with the resources of the mobile computing device. When docked at a desk top station, the screen of a mobile computing device is small, located at a relatively far distance from an operator and positioned such that an operator often has to place his or her back, neck or head in an uncomfortable position to adequately see the screen.\nErgonomic studies of computer-human operator interfaces have determined that a preferred monitor or screen height is that at which the top of the screen is even with the horizontal line of sight of the operator. The preferred angle of the screen is at 90 degrees to the line of sight.\nOne attempt to alleviate the undesirably low screen height of docked mobile computing devices includes providing a regular desk top monitor on a stand above the docking station. Though this may alleviate some of the problems of low-level mobile computing device screen height, the additional monitor and stand are often undesirably expensive and consumptive of space. This arrangement may also provide insufficient adjustment of screen viewing angle."}
{"text": "Field of the Invention\nThe present invention relates to a conductive material and a substrate having a conductive film formed thereon by using the conductive material.\nDescription of the Related Art\nA polymer having a conjugated double bond (i.e. π-conjugated polymer) does not show a conductivity by itself; however, if an appropriate anionic molecule is doped therein, it can express a conductivity, thereby giving a conductive polymer material (i.e. conductive polymer composition). As to the π-conjugated polymer, polyacetylene, (hetero) aromatic polymers such as polythiophene, polyselenophene, polytellurophene, polypyrrole, and polyaniline; a mixture thereof, etc., are used; and as to the anionic molecule (dopant), an anion of sulfonic acid type is most commonly used. This is because a sulfonic acid, which is a strong acid, can efficiently interact with the aforementioned π-conjugated polymers.\nAs to the anionic dopant of sulfonic acid type, sulfonic acid polymers such as polyvinyl sulfonic acid and polystyrene sulfonic acid (PSS) are widely used (Patent Document 1). The sulfonic acid polymer includes a vinylperfluoroalkyl ether sulfonic acid typified by Nafion (registered trademark), which is used for a fuel cell.\nPolystyrene sulfonic acid (PSS) has a sulfonic acid as a repeated monomer unit in the polymer main chain, so that it has a high doping effect to the π-conjugated polymer, and also can enhance water dispersibility of the π-conjugated polymer after being doped. This is because the hydrophilicity is kept due to the sulfo groups excessively present in PSS, and the dispersibility into water is therefore enhanced dramatically.\nPolythiophene having PSS as a dopant exhibits high conductivity and can be handled as an aqueous dispersion, so that it is expected to be used as a coating-type conductive film material in place of ITO (indium-tin oxide). As mentioned above, however, PSS is a water-soluble resin, and is hardly soluble in an organic solvent. Accordingly, the polythiophene having PSS as a dopant also has a high hydrophilicity, but a low affinity to an organic solvent and an organic substrate, and thus, it is difficult to disperse it into an organic solvent or to form a film onto an organic substrate.\nBesides, when the polythiophene having PSS as a dopant is used in, for example, a conductive film for an organic EL lighting, a large quantity of water tends to remain in the conductive film and the conductive film thus formed tends to absorb moisture from an outside atmosphere since the polythiophene having PSS as a dopant has an extremely high hydrophilicity as mentioned above. As a result, the problems arise that the luminous body of the organic EL chemically changes, thereby the light emitting capability is deteriorated, and that water agglomerates over time and defects are caused, which results in shortening of the lifetime of the whole organic EL device. Furthermore, there arise other problems in the polythiophene having PSS as a dopant that particles in the aqueous dispersion becomes large, the film surface becomes rough after the film formation, and a non-light emitting region, called dark spot, is caused when used for the organic EL lighting.\nIn addition, since the polythiophene having PSS as a dopant has an absorption at a wavelength of about 500 nm in the blue region, in the case that this material is used as a film coating a transparent substrate such as a transparent electrode, there arises another problem that when the conductivity required for the device to function is made up by the solid concentration or the thickness of the film, transmittance of the film is affected.\nPatent Document 2 discloses a conductive polymer composition composed of a conductive polymer which contains a π-conjugated polymer formed of a repeating unit selected from thiophene, selenophene, tellurophene, pyrrole, aniline, and a polycyclic aromatic compound, and a fluorinated acid polymer which can be wetted by an organic solvent and 50% or more of which is neutralized by a cation; and it is shown that an aqueous dispersion of the conductive polymer can be obtained by combining water, a precursor monomer of the π-conjugated polymer, the fluorinated acid polymer, and an oxidant, in any order.\nHowever, in such a conventional conductive polymer, particles are agglomerated in the dispersion immediately after synthesis. Also, if an organic solvent served as a conductive enhancer is added thereto to give a coating material, the agglomeration is further facilitated, so that the filterability thereof is deteriorated. If the coating material is applied by spin coating without filtration, a flat film cannot be obtained due to the effect of the particle agglomeration; and as a result, the problem of coating defect is caused.\nMoreover, development has been promoted in a flexible device. As a transparent conductive film for the current hard devices, ITO is widely used. ITO is, however, a crystalline film, and therefore there arises cracks in case of bending. Accordingly, it is a pressing need to develop a flexible transparent conductive film substituting for ITO. Polythiophene having PSS as a dopant forms a flexible film with high transparency, but involves a problem of low conductivity compared to ITO in addition to the aforementioned problem of dark spot.\nPatent Document 3 discloses a transparent conductive film using silver nanowires. The transparent conductive film using silver nanowires is one of a candidate for a conductive film for a flexible device, since it has high conductivity and transparency. The film using silver nanowire, however, conducts electricity only through the wire part, and therefore causes a problem that the light emission occurs at the wire parts only, not the whole surface when it is applied to an organic EL lighting."}
{"text": "Designing systems to test large numbers of semiconductor or electronic devices in a factory environment where each individual device under test generates significant thermal load is not a trivial task. Examples of such devices include but are not limited to high-speed microprocessors, laser diodes, and linear high-power amplifiers.\nFrom a semiconductor manufacturers standpoint, factory floor space often comes at a premium. A testing system that is capable of testing the largest number of semiconductor devices in the smallest factory floor footprint is therefore highly desirable. Removing large amounts of heat from high-dissipation devices in a very small footprint is problematic, however. Air-cooling devices undergoing testing require large heat sinks with a corresponding large surface area and size. Water-cooling devices undergoing testing require a lower corresponding heat exchanger area than air, but is still not the smallest solution possible and also has the intrinsic risk of damaging sensitive electronics in the test system should a leak develop.\nOne solution for cooling heat generating devices utilizes a dual-phase refrigerant where the physical process of evaporation (latent heat of vaporization) absorbs significantly more heat than air or fluid can alone for the same corresponding heat exchanger size. Cooling schemes that depend on either a refrigerant cycle or Carnot heat cycle have long been in existence. Various semiconductor test system manufacturers have used refrigerant cycles wrapped around gas compressors for their central plants.\nA key difference between a true refrigerant cycle and a Carnot heat cycle is nothing more than the thermodynamic direction that heat moves: low to high versus high to low. In a refrigerant cycle a compressor takes low pressure gas, compresses it generating high-pressure gas and heat, and then removes the heat in the condenser on the high-pressure side of the system. A Carnot heat cycle on the other hand, uses a fluid pump to pressurize liquid refrigerant into the evaporators. Heat then flashes the liquid refrigerant into gas and then a water sub-cooled condenser downstream condenses refrigerant gas on the low-pressure side of the system to complete the cycle. A classic example of a Carnot heat cycle is a fossil or nuclear-fueled steam power plant generating electricity.\nSome of the common problems experienced by current thermal systems used in semiconductor testing environments include refrigerant leaks, poor chiller mean-time-before-failure (MTBF) and inconsistent thermal performance as the device-under-test (DUT)/system thermal load parameters changes. Other problems associated with some thermal systems include the inability to charge the thermal system while the system is operational and the inability to control suction pressure to increase cooling temperature differential."}
{"text": "Modern hard disk drives may have one or more rotating disks. Each disk may have two disk surfaces in use. Hard disk drives also include a servo controller that drives a voice coil actuator to position a read-write head near a track on the surface of a rotating disk in order to read and/or write data to the track. The read-write head may be a magnetic transducer adapted to read data from and/or write data to the one or more rotating disks. Read operations are generally accomplished by sensing a magnetic field associated with the surface of the disk. Write operations are generally accomplished by selectively magnetizing the surface of the disk. The read-write head communicates with the servo controller, providing feedback, which is used in controlling the read-write head's positioning relative to the track. The read-write head is embedded in a slider, which flies on a thin layer of air, known as an air bearing, a very short distance above the rotating disk surface.\nThe trend in the hard disk drive industry is to increase the areal density of the rotating disk surfaces. This is usually achieved by decreasing the width of the tracks and reducing the flying height of the read-write head above the rotating disk surface, thus reducing the air bearing. Currently read-write heads fly at about 10 nanometers (nm) from the rotating disk surfaces."}
{"text": "The invention relates to voter ballots, which create machine-readable codes when a selection is registered for more accurate machine counting and to eliminate confusion of the voters\"\" intent.\nBriefly stated the invention is an election ballot comprising a plurality of political office candidate names printed in visible ink and a unique bar code printed for each candidate for a voter to enter a candidate selection; and a unique ballot number and bar code of that number printed on the ballot. A plurality of the ballots can be used together with a computer data base and an internet server for each voter to confirm his candidate selections were correct."}
{"text": "In a typical ink jet recording or printing system, ink droplets are ejected from a nozzle at high speed towards a recording element or medium to produce an image on the medium. The ink droplets, or recording liquid, generally comprise a recording agent, such as a dye or pigment, and a large amount of solvent. The solvent, or carrier liquid, typically is made up of water, an organic material such as a monohydric alcohol, a polyhydric alcohol or mixtures thereof.\nA continuing problem with ink jet printers is the accumulation of ink on ink jet nozzle plates, particularly around the orifice from which ink drops are ejected. The result of ink drops accumulating around the orifice is that it becomes wettable causing ink drops to be misdirected, degrading the quality of the printed image.\nTo limit or prevent the spreading of ink from the orifice to the nozzle plate, it is common practice to coat the ink jet nozzle plate with an anti-wetting layer. Examples of anti-wetting layers are coatings of hydrophobic polymer materials such as Teflon.RTM. and polyimide-siloxane, or a monomolecular layer of a material that chemically binds to the nozzle plate, e.g., alkyl thiols, alkyl trichlorosilanes and partially fluorinated alkyl silanes.\nInk jet nozzle plates are also contaminated by ink drops that land on the nozzle plate. These \"satellite\" ink drops are created as a by-product of the drop separation process of the primary ink drop that is used to print. Another source of contaminating ink are tiny ink drops that are created when the primary ink drop impacts recording material. Ink drops accumulating on the nozzle plate can also potentially attract contaminants such as paper fibers which cause the nozzles to become blocked. Partially or completely blocked nozzles can lead to missing or misdirected drops on the print media, either of which degrades the quality of the print.\nIn order to solve this problem, the nozzle plates are periodically wiped clean. Several wiping methods are known including wet wiping techniques utilizing inks as a cleaning solvent. While inks and ink solvents used to dilute inks may be used as a cleaning liquid, they are not optimized for this purpose. Inks may contain additives such as, for example, ethylene glycol, diethylene glycol, and diethylene glycol monobutyl ether which may be environmentally undesirable when released during cleaning in unventilated areas such as a home or an office.\nFurther, inks often contain various materials which may leave an undesirable residue on the ink jet print head nozzle plate. Thus while wiping removes ink drops from the nozzle plate, the hydrophobic anti-wetting coating on the nozzle plate may be severely contaminated and compromised by ink residue. The ink-fouled coating is therefore unable to prevent the spreading of ink from orifices.\nIt has also been discovered that hydrophobic coatings on an ink jet print head nozzle plate are susceptible to fouling by certain ink jet inks, such as those containing copper phthalocyanine dyes. The fouling of the nozzle plate by the ink can lead to excessive spreading by ink on to the nozzle plate during normal use, further aggravating drop placement problems. Another disadvantage in using inks as a cleaning solution is that they are expensive.\nThere remains a need for a simple, economical ink jet nozzle plate cleaning solution that will help maintain the anti-wetting character of ink jet nozzle plates so that an ink jet print head will consistently deliver accurate and reproducible drops of ink to a receiver resulting in photographic quality images."}
{"text": "The present invention relates to a semiconductor device and a method of manufacturing the same, more specifically, a semiconductor device comprising a nonvolatile memory of the stacked gate structure and a transistor of the single-layer gate structure, and a method of manufacturing the semiconductor device.\nThe logic semiconductor device combined with a nonvolatile semiconductor memory forms product fields, as of CPLD (Complex Programmable Logic Device) and FPGA (Field Programmable Gate Array), and because of their characteristic of programmability, so far have formed large markets because of their characteristic, programmability.\nThe logic semiconductor device combined with a nonvolatile memory has, in addition to flash memory cells, a high-voltage transistors for controlling the flash memory and low-voltage transistors of high-performance logic circuit integrated on the same semiconductor chip. The flash memory cells have gate electrode of the stacked structure of a floating gate and a control gate laid one on the other which is different from the single-layer structure of the high-voltage transistors and the low-voltage transistors. Accordingly, the process of manufacturing the logic semiconductor device combined with the nonvolatile memory requires the process specialized in forming the nonvolatile memory transistors of the stacked gate structure without changing characteristics of the peripheral circuits, especially the logic transistors.\nIn the usual combined process, the floating gates of the nonvolatile memory transistors are formed of the first-level conductive film (the first conductive film), and the control gates of the nonvolatile memory transistors and the gate electrodes of the peripheral transistors are formed of the second-level conductive film (the second conductive film). Then, the peripheral transistors are formed after the nonvolatile memory transistors have been formed, so as to prevent the process of manufacturing the nonvolatile memory transistors from influencing characteristics of the logic transistors. In terms of the process of forming the gate electrodes, after the second conductive film in the memory cell region have been patterned to form the control gates, the second conductive film in the peripheral circuit region is patterned to form the gate electrodes of the peripheral transistors.\nThe related arts are described in, e.g., Reference 1 (Japanese published unexamined patent application No. Hei 10-209390).\nHowever, the inventors of the present application have examined the process of fabricating the logic semiconductor device combined with the nonvolatile memory and found that the process causes the disadvantage that the second conductive film in the memory cell region is etched when the second conductive film is patterned to form the gate electrodes of the peripheral transistors."}
{"text": "Electronic devices, such as mobile computing devices, including tablets and smartphones, are typically encased in packaging at the place of manufacture or final assembly. Such packaging protects the devices during transportation and storage, while at the same time provides an attractive appearance to a consumer. However, consumer packaging is typically sealed, resulting in the electronic device being inaccessible from the time of packaging until the package is opened by the purchaser. As a result, the manufacturer is unable to access the sealed device once packaging has occurred. In some cases, it may be desirable to access the electronic device after packaging but before delivery to the customer.\nAccordingly, there is a need to provide improved servicing capability for packaged electronic devices."}
{"text": "This invention relates to a method for laminating a transparent safety panel to the screen-bearing viewing window of a CRT (cathode-ray tube) and particularly, but not exclusively, to a large CRT, and to the product of that method. By \"large CRT\" is meant a television picture tube or an information display tube having a viewing window bearing a viewing screen of at least a 30V size.\nIn one prior type of CRT, a glass safety panel is laminated to the viewing window of the CRT in order to reduce the danger of implosion and also, should the CRT implode, to reduce the danger of injury to people and things nearby. Suitable structures and methods for laminating CRTs smaller than 26V screen dimension have been described previously, for example, in U.S. Pat. No. 4,204,231 issued May 20, 1980 to M. M. Permenter.\nIn one prior laminating method, a safety panel is positioned in a desired spaced-apart relationship with a CRT window, and then a strip of flexible pressure-sensitive tape is wrapped around the edges of the CRT and panel to define a closed cell between the window and the panel. Thereafter, the cell is filled with a catalyzed liquid resin and allowed to cure to a clear transparent solid which adheres to the facing surfaces of the panel and the window. During the filling step, vent holes are punctured in the tape to allow air in the cell to escape. After the cell is filled with resin, the holes are taped shut to prevent both the leakage of resin and the formation of bubbles in the cell during the curing step. A foam tape with adhesive on both sides adhered to the margins of the panel and the window may replace the tape in the foregoing method.\nIn the foregoing method, the cell is filled with the viewing window positioned in a generally horizontal attitude, with the window facing downward. The tape provides a temporary hydraulic seal for the cell, and is also sufficiently strong to hold the panel temporarily in the desired downward-facing, spaced-apart relationship while the cell is being filled and the resin cures. In practice, especially when safety panels are laminated to CRT windows larger than about 25V-size, the taping step is not efficient and many temporary seals leak resin, and some seals fail to hold the safety panels in place. Also, because the windows face downward, it is difficult to determine whether gas bubbles are present in the viewable area in front of the window.\nBecause of the size and weight of a large CRT as defined above, all of these problems are aggravated and these prior methods are not practical for laminating a safety panel to the window of a large CRT. For example, a 25V-size CRT weighs about 55 pounds, while a 33V-size CRT weighs about 130 pounds and therefore cannot be handled manually in the factory. A cell, formed as described above, on a 33V-size CRT is difficult to fill with liquid resin with the window facing down because of the greater weight of the safety panel and because the greater weight of the resin causes greater leakage of resin during the filling and curing steps, especially through the venting holes in the tape. The tapes used in the prior methods to form the cell are not reliable to hold an 11-pound safety panel and about 8 pounds of liquid resin in the desired spaced apart relationship during the filling and curing steps. Sagging and wedging of the panel position and leakage and \"flow-out\" of the resin occur frequently with prior methods."}
{"text": "With the increasing public attentions to ecological problems, greening projects in areas such as deserts, slopes, embankments and wastelands are more and more concerned by the society. One of the typical greening methods available is a fabric-spreading method in which an ecological fabric is laid on the surface of the area to be afforested. In some conditions (e.g. soil slopes, etc.), the method is to secure the soil/sands using the ecological fabrics, and plant growth materials are placed under the ecological fabrics; and in other conditions (e.g. deserts, etc.), the method is to place the plant growth materials on the ecological fabrics. However, in any conditions, the greening is made by means of the ecological fabrics, therefore, it is critical in the whole greening process how to secure the ecological fabrics. In the traditional methods, the ecological fabric is fixed to the area to be afforested by an anchor rod piercing the ecological fabric, in this case, the anchor rod may be loosened due to the decrease in the friction force between the anchor rod and the sands or soil, and even be drawn out from the sands or soil, leading to the unsteady fixation of the ecological fabric, and hence to the displacement of the fabric. Such a method merely relying on the anchor rod has a poor effectiveness, and is also not suitable for zones such as deserts. It is also provided a method for fixing the ecological fabric by making trenches in desert areas, but it is needed to fill ecological bags into the trenches to press the ecological fabric. However, the production of the ecological bags is complicated and costly, which apparently does not meet the current needs and is not convenient to marketing promotion."}
{"text": "1. Field of the Invention\nThe present invention relates to a gas compressor for use in a vehicle air conditioning system or the like and, more particularly, aims to prevent fretting wear due to vibration of an oil separation filter contained in an apparatus to thereby achieve an improvement in durability.\n2. Description of the Related Art\nConventionally, a gas compressor of this type has, as shown, for example, in FIG. 7, a compression mechanism C in a compressor case 1, and this compression mechanism C sucks in low pressure refrigerant gas in a suction chamber 6 and compresses it before discharging it into a discharge chamber 5 as high pressure refrigerant gas. At this time, the high pressure refrigerant immediately after being discharged from the compression mechanism C contains oil in the mist state filled for the purpose, for example, of lubricating the sliding portion of the compression mechanism portion C, and the oil ingredient in such high pressure refrigerant gas is separated by an oil separator 20 contained in the gas compressor.\nThe oil separator 20 has an oil separation filter 21 formed by rolling a wire net into a cylinder, and the oil ingredient in the high pressure refrigerant gas is separated when a high speed gas flow of high pressure refrigerant gas strikes upon the wire net constituting the oil separation filter 21 or by the centrifugal force generated when high pressure refrigerant gas whirls inside the oil separation filter 21.\nHowever, in the oil separator 20 of the conventional construction described above, in which a pulsating high speed gas flow incessantly strikes upon the wire net constituting the oil separation filter 21, if the wire net is not supported and fixed in a satisfactory manner, a vibration such as twist is generated in the entire wire net, with the result that the constituent wire materials of the wire net (wires) rub each other, and the wire net and the support member in contact therewith rub each other, resulting in the generation of so-called fretting wear of the wire net to damage the wire net.\nThe present invention has been made in view of the above problem. It is an object of the present invention to provide a gas compressor in which fretting wear due to vibration of the oil separation filter is prevented to thereby achieve an improvement in durability.\nTo achieve the above object, there is provided, in accordance with the present invention, a gas compressor comprising a compression mechanism portion which sucks in, compresses, and discharges refrigerant gas, a discharge chamber for temporarily storing the refrigerant gas discharged from the compression mechanism portion, a discharge passage which communicates the compression mechanism portion with the discharge chamber and which guides the refrigerant gas discharged from the compression mechanism portion to the discharge chamber, an oil separation filter which is arranged in the discharge chamber at the downstream end of the discharge passage and which is formed by shaping a lattice-like member into a cylindrical configuration, pressing means for pressing the end portions of the oil separation filter to elastically deform the oil separation filter sufficiently to suppress vibration thereof during use of the gas compressor, and an oil sump formed at the bottom of the discharge chamber and adapted to store the oil separated by the oil separation filter.\nIn the present invention, the oil separation filter as a whole is firmly fixed in a pressed state by the pressing force when the pressing means presses the end portions of the oil separation filter.\nHere, as the pressing means, it is possible to adopt a construction comprising an abutment block portion against which the bottom end portion of the oil separation filter abuts, an outer peripheral block portion having a filter accommodating hole into which the oil separation filter is inserted, and a pressurizing plate put over the forward end portion of the oil separation filter and mounted and fixed to the outer peripheral block portion, wherein the forward end portion of the oil separation filter is provided so as to protrude from the filter accommodating hole on the surface of the outer peripheral block portion, and wherein the protruding forward end portion of the oil separation filter is pressurized and bent by the pressurizing plate to undergo elastic deformation.\nIn this construction, the forward end portion of the oil separation filter is pressurized and bent toward the abutment block portion by the pressurizing plate to undergo elastic deformation, and due to this elastic deformation of the oil separation filter, the oil separation filter is strongly pressed against the pressurizing plate, the abutment block, and the outer peripheral block portion and fixed thereto.\nIt is possible to adopt a construction in which the forward end portion of the oil separation filter is previously bent in a fixed direction by a bending amount smaller than a press-bending amount before being press-bent by the pressurizing plate.\nIn this construction, the forward end portion of the oil separation filter is press-bent by the pressurizing plate such that all the portions thereof are directed in the same direction, so that the construction proves superior in operability in fixing the bending direction.\nFurther, it is desirable that the direction in which the forward end portion of the oil separation filter is previously bent be toward the inner side of the filter.\nIn this construction, the forward end portion of the oil separation filter is press-bent by the pressurizing plate such that it is always directed to the inner side of the filter, so that there is no danger of the bent forward end portion of the oil separation filter being caught between the pressurizing plate and the outer peripheral block, making it possible to correctly mount and fix the pressurizing plate to the outer peripheral block portion.\nIt is possible to provide in the inner surface of the filter accommodating hole a clearance groove which allows entrance of the forward end portion of the oil separation filter when it is press-bent to the filter outer side.\nIn this construction, for example, there is no directivity in the bending of the forward end portion of the oil separation filter press-bent by the pressurizing plate, and even when the forward end portion of the oil separation filter is bent so as to be directed to the filter outer side, this press-bent portion enters the clearance groove, so that there is no danger of the bent forward end portion of the oil separation filter being caught between the pressurizing plate and the outer peripheral block, making it possible to correctly mount and fix the pressurizing plate to the outer peripheral block portion.\nThe oil separation filter may, for example, be formed of a lattice-like member, such as wire net or expanded metal; in particular, in the case of a wire net type oil separation filter, it is desirable to fuse by spot welding the regions where the constituent wire materials of the wire net (wires) intersect each other in order to prevent fretting wear as a result of the constituent wires of the wire net being rubbed against each other.\nFurther, the spot welding is performed in order to prevent fretting wear mainly due to vibration between the constituent wire materials of the wire net, and, from the viewpoint of the object of this spot welding, it is desirable to increase the number of spot welding points for the region where, for example, the high speed gas flow of high pressure refrigerant gas first strikes directly."}
{"text": "Speaker recognition generally includes the tasks of speaker identification and speaker verification: speaker identification involves labeling an unknown voice as one from a set of known voices, while speaker verification involves determining whether an unknown voice matches the known voice of a speaker whose identity is being claimed. In particular, speaker identity verification based on a person's voice is of considerable interest for providing telephone access to such services as banking transactions, credit card verification, remote access to dial-up computer databases.\nThe specific problem for which the invention has application is providing verified access to medical records databases for health care providers. For example, a voice log-in system would be used by doctors to obtain verified access to a number of medical databases at different health care institutions.\nTwo approaches to voice verification technology have been used: voice/mechanical and total voice. In a voice/mechanical system, a user manually enters a digit code and then speaks a voice password--the voice password is compared with a reference acoustic recognition model retrieved based on the secret code. In a total voice system, a user speaks a digit code--speaker-independent digit recognition is used to recognize the digits, and then a reference acoustic recognition model is loaded based on the recognized digit string, and voice verification is performed.\nThe total voice system can be described as a template-based, fixed-text speaker verification task. It includes the subtasks: (a) enrollment--creating a reference template composed of a set of speech features, as a function of time, associated with the fixed-text digit code for each valid user; (b) verification--selecting an appropriate reference template based on speaker-independent recognition of the digit code, and comparing input speech with reference templates at equivalent points in time, with a decision being based on similarity between input and reference, integrated over time; and (c) reference update--adapting the reference templates to accommodate changes in a valid user's speech after successful verification.\nThe enrollment procedure builds an initial reference template for a speaker by capturing the digit code, typically by examining the speech energy profile and determining the end points in time. For each successful verification, the dynamic reference updating procedure averages the feature sets associated with the reference template and the input speech, and the reference template is updated accordingly, thereby accommodating changes in a persons speech input as familiarity with the system is achieved.\nBoth the voice/mechanical and the total voice speaker verification systems are disadvantageous in that each requires the use of a specific digit code (either entered or spoken). In the medical records database application, such a requirement is problematic to the point of impracticality. Doctors and other health care providers typically see patients at a number of hospitals, each of which has its own medical records database system. These systems typically assign each user a unique identification code, requiring each health care provider to remember a different identification code for each institution and database system.\nAccordingly, a need exists for a speaker verification system that is not based on the spoken input of an identification code."}
{"text": "(1) Field of the Invention\nThe present invention relates to a filing system including a data processing apparatus which captures image data, stores it on a storage device, transmits it to another data processing device, and copies it onto a copy sheet. The present invention also relates to a data processing method and a computer readable medium which are incorporated into the data processing apparatus in the filing system.\n(2) Description of the Related Art\nRecently, a number of centralized mainframe computers are shared by all users in an organization, and supplemented by workstations and personal computers located in departmental user rooms and private offices. With the increased number of machines has come the need to move data and files from one machine to another. A filing system is one approach to solving the file transfer problem. In the filing system, shared data is placed on a file server, and, when needed, individual machines are made to access data files located on the remote file server. This approach works well when the number of files that needs to be exchanged is not small.\nIn a conventional filing system, a data processing device (for example, a personal computer) is provided with a scanner, and the scanner captures image data by optically scanning the document to be copied or transmitted. However, when the number of documents that need to be exchanged between authorized users is large, the data capturing using the scanner becomes a considerably time-consuming task, and the access to the image data on the remote file server from the data processing device becomes a burdensome operation. This causes the delay of implementation of the conventional filing system. In addition, the contents of documents having confidential data, such as a password for the file access, may be revealed to the operator of the data processing device during the data capturing. The conventional filing system has such a problem of data security.\nBy taking account of the above-described problems, the inventors of the present invention have noted that image data of documents, which need to be stored in a filing system, are always first captured using scanners of copiers or facsimiles for corporate meeting uses or business purposes, and subsequently translated into usable signals using analog-to-digital converters. Hence, in order to eliminate the above-described problems, the inventors have designed a data processing apparatus and method, for use in the filing system, which creates identifications of the authorized users for such documents when the image data are first captured, and has the image data and the identifications stored on a storage medium of the filing system in a such a manner that the stored image data are subsequently accessible by the authorized users using the identifications. Further, the filing system must be designed to ensure security of the image data stored therein."}
{"text": "A type of material handling units useful for transferring workpieces from a chucking position to a loading position for a machine has a swinging arm or manipulator provided at its outer end with means for gripping the workpiece. In a material handling unit of this kind relatively high positioning accuracy is required at the chucking position where the swing motion starts and at the loading position where it ends.\nSuch positioning is commonly carried out by engagement of a striking face of the swinger arm against a fixed abutment mounted on the frame of the material handling unit. In this manner high positioning accuracy is attained in comparison with a mechanical cam or electric positioning device and moreover repeatable accuracy and reliability are extremely high.\nHowever when positioning the swinger arm by engagement with an abutment especially when the speed of the swinging motion of the swing arm is increased and the cycle time is reduced the swing arm may bounce or vibrate by reason of the impact shock in accordance with the weight of the payload. When such bouncing or vibration occurs the desired high positional accuracy and reliable motion cannot be attained."}
{"text": "1. Field of the Invention\nThe present invention relates generally to transmission of radio frequency (RF) signal waves and specifically to the combination of different RF signals into a common transmission line, such as a waveguide, for application to a common propagating device such as an antenna for further transmission of the electromagnetic wave.\n2. Background and Related Art\nIn electronic apparatus concerned with amplification and transmission of electromagnetic energy, the energy is usually transferred from one component of the system to another in the form of current through electrical conductors such as wire or cable. At very high frequencies and high power, the transfer of electrical energy through current carrying conductors becomes inefficient and impractical. Instead of transferring electrical energy as current through conductors, the high frequency high power energy can be transmitted by electromagnetic waves constrained within a guiding structure or waveguide. The transmission of energy by waveguide is carried out by radiating the energy, in the form of electromagnetic waves, into the inside of the waveguide structure through various coupling mechanisms such as probes, loops, crossbars and the like. This same energy may be extracted from the waveguide by a similar coupling mechanism at another location on the waveguide. Where the energy is in the form of current, conveyed in multiple conductors, devices such as coaxial cross-coupled hybrids may be used.\nMany applications exist where it is desired to combine signals of various frequencies into a common transmission line structure. One example of such an application is in the field of microwave communications. One requirement of such combination is that the different signals must not improperly interfere with each other.\nIn this regard, various tuned filtering devices have been used in the prior art. However, due to the non-ideal nature of some of such devices, unwanted distortions and losses are inevitably introduced into the frequency bands of the different signals, causing the signals to be degraded in quality.\nOne particular application for combining electromagnetic signal waves has arisen with the advent of digital television (DTV). Broadcasters who wish to provide DTV service still must provide broadcast signals in the conventional NTSC (National Television Systems Committee) format to accommodate viewers with conventional NTSC television receivers. Absent a feasible and efficient way to combine the DTV signal and the NTSC signal for transmission by a common broadcast antenna (possibly on adjacent frequency channels) without causing unacceptable interference, interaction and distortion of the signals, this will require some broadcasters to build additional antennae and towers to carry their DTV signals. This represents a very expensive and undesirable proposition.\nUnder present NTSC broadcast allocations, each active channel of 6 MHz bandwidth is separated from other channels by one or more 6 MHz wide inactive channels, in order to prevent interchannel interference between signals. Because of the limited nature of the electromagnetic spectrum and the requirement for broadcasters to continue to provide conventional NTSC signals to serve viewers who do not own DTV receivers, DTV signals will, in some cases, be broadcast on currently inactive 6 MHz wide channels adjacent to the upper side of current NTSC broadcasting channels.\nVarious techniques have been proposed to provide an acceptable combination of analog (i.e., NTSC) and digital DTV signals on adjacent broadcast channel allocations to meet this requirement. Proposed combiners utilizing wideband and other LC filters unavoidably introduce large signal phase delay and amplitude artifacts which are at best difficult to correct in order to reestablish the original signal purity specifications of the transmitter. The difficulty arises from the fact that a broadband filter is required for each signal having a \"cliff\" response directly at the common band edge between the combined signals, because there is virtually no guard band between the two signals. Such LC filters have much difficulty achieving good performance characteristics under such constraints, and the need for stability over varying environmental conditions (mainly temperature) adds further to the difficulty in designing acceptable combiners using these LC filters.\nAccordingly, there exists a need in the art to enable the combination of signals of different frequencies on adjacent broadcast channels without mutual interference and distortion."}
{"text": "1. Field of the Invention\nThe present invention relates to a shock absorbing type steering shaft which is capable of contracting in the axial direction upon collision of the car owing to the impact of the collision, and a method of manufacturing such steering shaft. Particularly, the present invention relates to a shock absorbing type steering shaft capable of stabilizing a collapse load which acts at the time of collision of the car even when the inner shaft of the car is a solid shaft so as to reduce the manufacturing cost, as well as a method of manufacturing such steering shaft.\n2. Related Background Art\nIn a steering apparatus for a car, there is provided a shock absorbing type steering shaft which can be contracted in the axial direction upon collision, by the impact of the collision.\nAn example of a method for manufacturing such shock absorbing type steering shaft is disclosed in Japanese Patent Application Laid-Open No. 52-25330, in which a solid inner shaft having a non-circular cross section in a substantially oval form is inserted into a tubular outer shaft having a non-circular cross section in a substantially oval form, the outer shaft is pressed by a pressing member inwardly from the outer side thereof in the radial direction to plastically deform the both shafts locally, and then both the shafts are pressure-fixed to each other elastically by the portions thus plastically deformed. In this manner, while a relative movement between the both shafts in the axial direction is prevented normally, the plastically deformed portions are caused to collapse at the time of collision of the car, whereby the inner shaft is moved into the outer shaft to contract the entire length of the steering shaft, thereby absorbing the shock.\nHowever, according to the manufacturing method disclosed above, since the plastic deforming is conducted in the state that the solid inner shaft is inserted in the outer shaft, an amount of plastic deformation of the outer shaft is not sufficient, so that there is a possibility of backlash in the pressure contact and fixation of both the shafts.\nIn Japanese Patent Application Laid-Open No. 1-58373, in the state that a jig is inserted in a tubular outer shaft having a circular cross section, the outer shaft is pressed inwardly in the radial direction by the pressing member from the outer side thereof, so that the outer shaft is plastically deformed locally to have a non-circular cross section. Subsequently, the jig is removed, and the solid inner shaft having a circular cross section is inserted into the outer shaft, whereby the inner shaft is pressure-fixed to the portion of the outer shaft which has been plastically deformed to the noncircular cross section.\nHowever, according to the manufacturing method disclosed above, it is required to once insert the jig into the outer shaft and, after the plastic processing, to remove the jig. Thus, the manufacturing process becomes complicated for such steps of inserting and removing the jig and, as a result, the manufacturing cost is increased.\nFurther, in Japanese Patent Application Laid-Open No. 9-272447, an inward protrusion is formed in the radial direction on a tubular outer shaft having a circular cross section with a female serration, an annular recess from which a serration has been removed is formed on the outer peripheral surface of a solid inner shaft having a circular cross section with a male serration, and the annular recess of the inner shaft is caused to engage with the inward protrusion in the radial direction of this outer shaft in such a manner that the both shafts are pressure-contacted and fixed with pressure to each other.\nHowever, according to the manufacturing method disclosed above, it is required to remove a part of the serration from the male serration formed on the outer peripheral surface of the inner shaft and to form the annular recess, which may result in an increase in the manufacturing cost.\nFurther, according to Japanese Patent Application Laid-Open No. 10-181615 (which has been converted from an application for a Japanese Utility Model Registration filed in 1992), the tip end of a tubular outer shaft having a circular cross section and the tip end of a solid inner shaft having a circular cross section are plastically deformed to have non-circular (flat or elliptic) cross sections, respectively. On the other hand, inner portions other than the tip ends of the both shafts are maintained to have circular cross sections. Then, the inner shaft is inserted into the outer shaft, and the portion having the non-circular cross section of the outer shaft is strongly pressure-contacted and fixed to the portion having the circular cross section of the inner shaft and the portion having the circular cross section of the outer shaft is strongly pressure-contacted and fixed to the portion having the non-circular cross section of the inner shaft, respectively, whereby the both shafts are pressure-contacted and fixed to each other.\nHowever, according to the manufacturing method disclosed above, the plastic deforming of the tip end of the outer shaft and the plastic deforming of the tip end of the inner shaft are conducted separately, so that the manufacturing process becomes complicated, which may bring about an increase in the manufacturing cost.\nFurther, in Japanese Patent Application Laid-Open No. 8-91230, like in Japanese Patent Application Laid-Open No. 10-181615, the tip end of the outer shaft and the tip end of the inner shaft are plastically deformed to have non-circular cross sections. However, as a step of this plastic deforming, the tip end of the inner shaft is inserted in the tip end of the outer shaft to be maintained in an overlapping state, and these both overlapping tip ends are pressed inwardly in the radial direction by the pressing members to be plastically deformed to have non-circular cross sections. Thus, the both tip ends are plastically deformed at a time to simplify the manufacturing process.\nHowever, according to the manufacturing method disclosed above, the inner shaft is required to be hollow and, when the inner shaft is solid, even if the tip end of the solid inner shaft is inserted into the tip end of the tubular outer shaft and the both overlapping tip ends are pressed inwardly in the radial direction by the pressing members, the plastic deforming can not be conducted satisfactorily with respect to the solid inner shaft.\nFurther, according to Japanese Patent Application Laid-Open No. 10-147245, like in Japanese Patent Application Laid-Open No. 8-91230 mentioned above, the both tip ends are plastically deformed by the pressing members at a time in the state that the both tip ends of the outer shaft and the inner shaft overlap each other. It is further arranged to adjust the collapse load by adjusting the press load of the pressing members.\nHowever, according to the manufacturing method disclosed above, like in Japanese Patent Application Laid-Open No. 8-91230, when the inner shaft is solid, the plastic deforming can not be conducted satisfactorily.\nIncidentally, when the inner shaft is a solid shaft, it is comparatively difficult to plastically deform the both tip ends at a time in the state that the tip ends of the both shafts overlap each other, like Japanese Patent Application Laid-Open No. 8-91230. However, a part of the above Japanese Patent Application Laid-Open No. 8-91230 discloses a method for pressure-contacting and fixing the solid inner shaft and the tubular outer shaft to each other, as shown in FIG. 19.\nA bottomed hole 103 is formed at the tip end 101a of the solid inner shaft 101 having a circular cross section with a male serration 102, and the tip end 101a of this inner shaft 101 is fitted in the tip end 104a of the tubular outer shaft 104 having a circular cross section with a female serration 105. These both overlapping tip ends 101a and 104a are pressed by the pressing members 106 and 106 from the upper and lower directions thereof so that the both tip ends 101a and 104a are plastically deformed to respectively have noncircular (flat or elliptic) cross sections, meanwhile the inner portions other than the tip ends of the both shafts 101 and 104 are maintained to have the circular cross sections.\nAfter that, the pressing members 106 and 106 are removed and the inner shaft 101 is further inserted into the outer shaft 104, the tip end 101a of the non-circular cross section of the inner shaft 101 is strongly pressure-contacted and fixed to the inner portion of the circular cross section of the outer shaft 104, and the tip end 104a of the non-circular cross section of the outer shaft 104 is strongly pressure-contacted and fixed to the inner portion of the circular cross section of the inner shaft 101.\nHowever, the drilling and processing works are required to form the bottomed hole 103 on the tip end 101a of the solid inner shaft 101, which may bring about an increase in the manufacturing cost.\nWhile the outer shaft 104 is obtained by plastically deforming a tubular shaft, the inner shaft 101 is obtained by plastically deforming a solid shaft on which the bottomed hole 103 is drilled. Thus, the level of the plastic deforming is not always the same for the both shafts, and one of the shafts may be deformed excessively. As a result, the force of pressure-contact and fixation is different between the both shafts, so that it becomes difficult to stabilize the collapse load to act upon collision of the car.\nThe present invention has been conceived taking the circumstances mentioned above into consideration, and an object of the invention is to provide a shock absorbing type steering shaft which can stabilize a collapse load acting upon collision of the car and can reduce the manufacturing cost, as well as a method for manufacturing such steering shaft.\nIn order to achieve the above object, there is provided according to the present invention, a shock absorbing type steering shaft comprising a tubular outer shaft having a female serration and a solid inner shaft having a male serration, characterized in that:\nafter the outer diameter surface of the outer shaft is pressed inwardly in the radial direction to plastically deform an axially inner portion thereof, the inner shaft is inserted into the outer shaft so that the inner shaft is pressure-contacted and fixed to the plastically deformed portion of the outer shaft.\nIn order to achieve the above object, there is also provided according to the present invention a method for manufacturing a shock absorbing type steering shaft for fitting and fixing a solid inner shaft having a male serration in a tubular outer shaft having a female serration, comprising the steps of:\ninserting the tip end of the inner shaft into the tip end of the outer shaft to hold the same in a state that these tip ends are overlapping each other;\npressing the outer shaft at which the both tip ends are not overlapping each other inwardly in the radial direction by means of a pressing member to plastically deform an axially inner portion thereof; and\nremoving this pressing member to further insert the inner shaft into the outer shaft, thereby pressure-contacting and fixing the inner shaft to the plastically deformed portion of the outer shaft.\nAs describe above, according to the present invention, the tip end of the inner shaft is inserted into the tip end of the outer shaft to be maintained in an overlapping manner, the inner portion of the outer shaft is pressed inwardly in the radial direction by the pressing members and is plastically deformed, and the inner shaft is pressure-contacted to be fixed to this plastically deformed portion of the outer shaft. That is, the outer shaft is plastically deformed while preventing the outer tube from being excessively deformed, by plastically deforming the inner portion of the outer shaft in the state that the tip ends of the both shafts overlap each other slightly. Consequently, it is possible to stably maintain the force of pressure contact and fixation of the steering shaft in its accomplished state. It is also possible to stabilize a collapse load which acts at the time of collision of the car. In addition, since a drilling step is unnecessary, unlike in the conventional method, the manufacturing cost can be reduced.\nAccording to another aspect of the present invention, there is provided a method for manufacturing a shock absorbing type steering shaft in which a solid inner shaft having a male serration is pressure-fitted in and fixed to a tubular outer shaft having the female serration, comprising the steps of:\npressing and deforming a predetermined portion of the tip end of the outer shaft from the outside thereof in the radial direction to form the female serration portion; and\npressing and fitting the male serration portion of the inner shaft in this pressed and deformed portion to pressure-contact and fix the inner shaft thereto. There is also provided a shock absorbing type steering shaft which is manufactured by such method."}
{"text": "1. Field of the Invention\nThis invention relates generally to a poker styled card game, and more specifically to a poker styled card game wherein a player attempts to achieve a high ranking 5-card poker hand by utilizing four cards dealt to that player and four community cards.\n2. Description of the Related Art\nCurrently, one of the most popular poker games played is a variation of poker called Texas Hold Em. In Texas Hold Em, each player is dealt two cards, with the face of the cards being disposed toward the playing surface, referred to as “hole cards.” In a standard Texas Hold Em game, once each player has received their two hole cards, a round of betting occurs. In some games there are two compulsory bets at the commencement of the first round of play: the person immediately to the left of the dealer must place a bet equal to one half of the minimum predetermined bet, referred to as the “blind” and the player immediately to the left of the blind must place a bet equal to the minimum predetermined bet for that round, referred to as the “big blind.” Play proceeds around the table in a clockwise manner, with each player given the opportunity to either fold, call the bet of the big blind, or raise the big blind. Any player that has not folded is referred to as an “active player.” Once all bets have been called, three cards are dealt with the face of these cards being disposed opposite the playing surface (referred to as “the flop”) to establish the “community” cards (cards that each player can use to complete their poker hand). Another round of betting occurs, with each active player having the opportunity to fold, bet, raise a previous bet made in that round, or check. Once all bets have been called, a fourth community card (referred to as “the turn”) is added to the community with the face of this card being disposed opposite the playing surface. Another round of betting commences. Again, once all bets are called, a fifth and final card (referred to as “the river”) is dealt with its face being disposed opposite the playing surface. A final round of betting occurs. Once all bets have been called, the showdown occurs, wherein all active players reveal their hole cards to determine the winner of the game based upon the highest ranking 5 card poker hand created from the 5 community cards and that player's two hole cards. The winner typically receives the proceeds of the wagering, commonly referred to as “the pot.”\nThere have been countless variations of poker, and more specifically, Texas Hold Em, which itself appears to be derived from 7 card stud poker. For example, U.S. Pat. No. 5,531,448 (Moody) describes a variation of Texas Hold Em wherein players attempt to make the highest 5-card poker hand out of a set of hole cards and a set of community cards. In this variation, however, three hole cards are dealt to each player, and eight cards are dealt, some with their faces disposed opposite the playing surface, and some with their faces disposed toward the playing surface. These eight community cards are arranged in four groups of two cards with the community card layout being in the form of a directional compass with a pair of community cards at each point of the compass. The positioning of the pair of community cards on these respective points of the compass determine which, if any, of the cards in the respective pairs are disposed with their face toward or opposite the playing surface. Each player then can choose which pair of community cards he or she wished to utilize in completing their 5-card poker hand.\nU.S. Pat. No. 5,489,101 (Moody) again describes a variation of Texas Hold Em, again with each player attempting to make the highest 5-card poker hand out of a set of hole cards and a set of community cards. In this variation each player is dealt 5 hole cards, with six cards being dealt with their face disposed toward the playing surface as the community. The community cards are arranged in a three row triangle pattern. A player may discard up to 5 of his or her hole cards. The dealer then positions the community cards with their faces disposed opposite the playing surface. Depending on how many hole cards a player discarded, predetermined cards in the community are used by that player to complete his or her 5-card poker hand.\nU.S. Pat. No. 5,718,430 (Aramapakul, et al.) discloses a variation of Texas Hold Em wherein each player is dealt 3 hole cards. The player must discard one hole card. Five cards are then dealt with their faces disposed opposite the playing surface to serve as the community. These two hole cards and five community cards are then used to create a 5-card poker hand.\nU.S. Pat. No. 6,042,188 (Poitra) discloses a variation wherein each player makes an initial wager prior to being dealt any hole cards. After this initial round of betting, each player is dealt two hole cards. Three cards are dealt, two with their faces disposed opposite the playing surface, and the remaining community card being disposed with its face disposed toward the playing surface, as the community cards. At this point, each player is given the opportunity to “double down.” Following this opportunity, the showdown occurs and all bets are resolved in favor of the player having the highest 5-card poker hand.\nAdditionally, U.S. Pat. No. 6,651,983 (Chobanian) discloses a variation wherein each player initially is dealt 2 hole cards. A round of betting then occurs. Four community cards are then dealt with their faces disposed opposite the playing surface. Another round of betting occurs. Finally, each player is dealt one final hole card. A final round of betting takes place. Once all bets have been called, the showdown occurs, with the winner being that active player with the highest 5-card poker hand.\nAmong the deficiencies in the prior art is that the variations of Texas Hold Em involve complex instructions, which can intimidate the novice player. Additionally, a deficiency in the prior art is the limited number of hole cards utilized. When the number of hole cards is limited to two or three as seen in Texas Hold Em and its variations, the strategy of the game is limited, as a player has only to guess two cards of his or her opponent.\nThus, there is a need in the relevant art for a modified game of poker that provides a new variation of Texas Hold Em allowing for the inclusion of plurality of players that is easy to learn, is fast paced, and increases the skill necessary in determining the relative strength of a player's hand versus that of his or her opponent's."}
{"text": "1. Field of the Invention\nThe present invention relates to a hanger including a body member and slide member slideably mounted on the body member and hanging articles such as bicycles or kayaks for display.\n2. Description of the Related Art\nA conventional hanger device as shown in FIGS. 11 and 12 includes a hanger 100, two slots 101 formed on two ends of the hanger 100 and two hooks 102 slideably disposed in the slots 101 respectively via two fasteners 103. Each hook 102 has a through-hole, and each fastener 103 is adapted to insert through the through-hole of the related hook 102 and to be engaged with the related slot 101. The fasteners 103 are adapted to tighten/loosen the hooks 102 with respect to the hanger 100 and allow fixing/sliding of the hooks 102 with respect to the slots 101. A clamping device 104 is disposed on the hanger 100 between the two ends of the hanger 100 and attaches the hanger 100 to a hanger stand (not shown).\nHowever, the structure strength of the hanger 100 could be weakened via the slots 101 which are formed on the hanger 100, and the hooks 102 are limited to slide along the related slot 101 only. The large size of each slot 101, the more weakened structure of the hanger 100. Moreover, each hook 102 is attached to the hanger 100 via the fastener 103 so that the hooks 102 can not slide along the related slot 101 in a smooth manner.\nThe present invention is therefore intended to obviate or at least alleviate the problems encountered in the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to a vibratory barrel finishing apparatus and, more particularly, to a long-travel annular vibratory barrel type finishing apparatus for line-processing.\nVibratory barrel finishing apparatuses are broadly divisible into two types: a box-type and a circular-type. The vibratory barrel finishing apparatus of the present invention possesses the advantages inherent in both these two types of apparatus, i.e. an oval path having adaptability for long line-finishing which is an advantage peculiar to the box-type apparatus and a spiral flow of mass which is an advantage of the circular-type apparatus.\n2. Description of the Prior Art\nGerman Patents DBP1047993, DBP1036092 and U.S. Pat. No. 3,071,900, disclose vibratory finishing apparatuses in which the line of the conventional box-type vibratory barrel is extended. In this known type of apparatus, however, a huge driving means is required for effecting the driving, separation and circulation of mass, resulting in an impractically high cost of production. A circular type apparatus having an extended line is disclosed in Japanese Patent Publication No. 16558/1975, as well as in Japanese Patent Laid-open Publication No. 141995/1978. The length of the extended line, however, is only twice as large as that of the conventional barrel.\nThe vibratory barrel finishing apparatus disclosed in the Japanese Patent Publication No. 16558/1975 has a vibrator disposed at the center of the vibratory machine body. The finishing barrel has a corner angle which is a right angle or an acute angle and has a separating device projecting from the barrel body. Thus, this apparatus has an assymmetrical construction and, therefore, cannot be incorporated in a processing line. It is presumed that this apparatus cannot provide a satisfactory processing effect.\nThe Japanese Patent Laid-open Publication No. 141995/1978 states that \"It makes possible to incorporate not single vibrator but a plurality of vibrator in the space extending along the length of the barrel.\" Judging from the above-statement, as well as from the attached FIGS. 1-9 and associated description which lacks a consideration of synchronized rotation of unbalanced weights, it is considered that the inventor of this invention has made a serious mistake. Namely, it is known that non-synchronized vibrations of a plurality of vibrators applied to the mass in the barrel causes a random movement of the mass or overflowing of the mass from the barrel. The aforementioned Japanese Patent Laid-open Publication No. 141995/1978 teaches how to overcome the above-mentioned problem by inclining the barrel wall outwardly. This, however, complicates the construction extraordinarily and causes a rise of the cost, as well as suppression of the finishing effect due to an obstruction of the flow of the mass."}
{"text": "Many medical procedures involve the extraction and replacement of flowing blood from, and back into, a donor or patient. The reasons for doing this vary, but generally, they involve subjecting the blood to some process that cannot be carried out inside the body. When the blood is outside the patient it is conducted through machinery that processes the blood. The various processes include, but are not limited to, hemodialysis, hemofiltration, hemodiafiltration, blood and blood component collection, plasmaphresis, aphresis, and blood oxygenation.\nOne technique for extracorporeal blood processing employs a single “access,” for example a single needle in the vein of the patient or a fistula. A volume of blood is cyclically drawn through the access at one time, processed, and then returned through the same access at another time. Single access systems are uncommon because they limit the rate of processing to half the capacity permitted by the access. As a result, two-access systems, in which blood is drawn from a first access, called an arterial access, and returned through a second access, called a venous access, are much faster and more common. These accesses include catheters, catheters with subcutaneous ports, fistulas, and grafts.\nThe processes listed above, and others, often involve the movement of large amounts of blood at a very high rate. For example, 500 ml. of blood may be drawn out and replaced every minute, which is about 5% of the patient's entire supply. If a leak occurs in such a system, the patient could be drained of enough blood in a few minutes to cause loss of consciousness with death following soon thereafter. As a result, such extracorporeal blood circuits are normally used in very safe environments, such as hospitals and treatment centers, and attended by highly trained technicians and doctors nearby. Even with close supervision, a number of deaths occur in the United States every year due to undue blood loss from leaks.\nLeaks present a very real risk. Leaks can occur for various reasons, among them: extraction of a needle, disconnection of a luer, poor manufacture of components, cuts in tubing, and leaks in a catheter. However, in terms of current technology, the most reliable solution to this risk, that of direct and constant trained supervision in a safe environment, has an enormous negative impact on the lifestyles of patients who require frequent treatment and on labor requirements of the institutions performing such therapies. Thus, there is a perennial need in the art for ultra-safe systems that can be used in a non-clinical setting and/or without the need for highly trained and expensive staff. Currently, there is great interest in ways of providing systems for patients to use at home. One of the risks for such systems is the danger of leaks. As a result, a number of companies have dedicated resources to the solution of the problem of leak detection.\nIn single-access systems, loss of blood through the patient access and blood circuit can be indirectly detected by detecting the infiltration of air during the draw cycle. Air is typically detected using an ultrasonic air detector on the tubing line, which detects air bubbles in the blood. The detection of air bubbles triggers the system to halt the pump and clamp the line to prevent air bubbles from being injected into the patient. Examples of such systems are described in U.S. Pat. Nos. 3,985,134, 4,614,590, and 5,120,303.\nWhile detection of air infiltration is a reliable technique for detecting leaks in single access systems, the more attractive two-access systems, in which blood is drawn continuously from one access and returned continuously through another, present problems. While a disconnection or leak in the draw line can be sensed by detecting air infiltration, just as with the single needle system, a leak in the return line cannot be so detected. This problem has been addressed in a number of different ways, some of which are generally accepted in the industry.\nThe first level of protection against return line blood loss is the use of locking luers on all connections, as described in International Standard ISO 594-2 which help to minimize the possibility of spontaneous disconnection during treatment. Care in the connection and taping of lines to the patient's bodies is also a known strategy for minimizing this risk.\nA higher level of protection is the provision of venous pressure monitoring, which detects a precipitous decrease in the venous line pressure. This technique is outlined in International Standard IEC 60601-2-16. This approach, although providing some additional protection, is not very robust, because most of the pressure loss in the venous line is in the needle used to access the patient. There is very little pressure change in the venous return line that can be detected in the event of a disconnection, so long as the needle remains attached to the return line. Thus, the pressure signal is very weak. The signal is no stronger for small leaks in the return line, where the pressure changes are too small to be detected with any reliability. One way to compensate for the low pressure signal is to make the system more sensitive, as described in U.S. Pat. No. 6,221,040, but this strategy can cause many false positives. It is inevitable that the sensitivity of the system will have to be traded against the burden of monitoring false alarms. Inevitably this leads to compromises in safety. In addition, pressure sensing methods cannot be used at all for detecting small leaks.\nYet another approach, described for example in PCT application US98/19266, is to place fluid detectors near the patient's access and/or on the floor under the patient. The system responds only after blood has leaked and collected in the vicinity of a fluid detector. A misplaced detector can defeat such a system and the path of a leak cannot be reliably predicted. For instance, a rivulet of blood may adhere to the patient's body and transfer blood to points remote from the detector. Even efforts to avoid this situation can be defeated by movement of the patient, deliberate or inadvertent (e.g., the unconscious movement of a sleeping patient).\nStill another device for detecting leaks is described in U.S. Pat. No. 6,044,691. According to the description, the circuit is checked for leaks prior to the treatment operation. For example, a heated fluid may be run through the circuit and its leakage detected by means of a thermistor. The weakness of this approach is immediately apparent: there is no assurance that the system's integrity will persist, throughout the treatment cycle, as confirmed by the pre-treatment test. Thus, this method also fails to address the entire risk.\nYet another device for checking for leaks in return lines is described in U.S. Pat. No. 6,090,048. In the disclosed system, a pressure signal is sensed at the access and used to infer its integrity. The pressure wave may be the patient's pulse or it may be artificially generated by the pump. This approach cannot detect small leaks and is not very sensitive unless powerful pressure waves are used; in which case the effect can produce considerable discomfort in the patient.\nClearly detection of leaks by prior art methods fails to reduce the risk of dangerous blood loss to an acceptable level. In general, the risk of leakage-related deaths increases with the decrease in medical staff per patient driven by the high cost of trained staff. Currently, with lower staffing levels comes the increased risk of unattended leaks. Thus, there has been, and continues to be, a need in the prior art for a foolproof approach to detection of a return line leak or disconnection.\nIn an area unrelated to leak detection, U.S. Pat. No. 6,177,049 B1 suggests the idea of reversing the direction of blood flow for purposes of patency testing and access-clearing. The patency tests alluded to by the '049 patent refer simply to the conventional idea of forcing blood through each access to clear occlusions and to ascertain the flow inside a fistula."}
{"text": "1. Field of the Invention\nThis invention relates to the art of dispensing pressure sensitive labels.\n2. Brief Description of the Prior Art\nVarious U.S. Pat. Nos. 1,642,387, 2,259,358, 2,275,064, 2,516,487, 2,620,205, 3,051,353, 3,265,553, 3,501,365, 3,551,251, and 3,611,929 and British Pat. No. 1,057,126, Feb. 1, 1967 are made of record. U.S. Pat. No. 3,501,365 British Pat. No. 1,057,126 disclose U-shaped cuts in the supporting web or strip. In connection with one embodiment of U.S. Pat. No. 3,501,365 for example, the patent discloses \"that the feed holes are die cut in such a manner that the leading edge of the feed hole has been cut whereas the trailing edge has not\" been cut. The loose flap or flap portions formed by the U-shaped cut extends in a leading direction as the flap portion moves toward the peel edge or separator point. Due to the fact that the flap portion is adhered to the adhesive on the overlying label and due to the fact that the flap portions extend in the leading direction, the flap portions fold out of the plane of the supporting web at the peel edge. Also, the flap portions adhere to the overlying label so well because when the flap-forming cuts are made, the supporting web and the label are driven together by the cutting knives into intimate contact with the intervening adhesive. This increases the holding force of the flap portions to the overlying labels and consequently diminishes the tendency of the flap portions to separate from the overlying labels at the peel edge. This sometimes results in tearing at the ends of the U-shaped cut. As the supporting material at the ends of a U-shaped cut tears, the flap portion grows in area. As the flap portion grows, the enlarged flap portion is held to the overlying label or labels by an increased area of adhesive. Accordingly, the force holding the enlarged flap portion to the overlying label or labels increases as the tearing progresses until eventual rupture or breakage of the supporting web. In another embodiment of U.S. Pat. No. 3,501,365, loose internal portions formed by annular cuts move out of the plane of the supporting web at the peel edge along with the leading label as the supporting web is drawn about the peel edge. These loose internal portions cannot tear the supporting web because they are severed therefrom by the annular cut. However, in both embodiments, the flap portions move out of the plane of the supporting web at the peel edge for the same reasons."}
{"text": "Seat belts, also known as safety belts, are safety devices designed to secure and restrain the occupant of a vehicle, preventing harmful movement during a collision or a sudden stop. The seat belt reduces the likelihood of death or serious injury in traffic collisions by (a) reducing the force of secondary impacts between the occupant and interior surfaces of the vehicle, (b) preventing the occupant being ejected from the vehicle in a crash, and (c) keeping the occupant positioned correctly for maximum effectiveness of airbags where provided.\nSeat belts come in various types, including lap restraints or lap belts, sash belts or three point seatbelts combining the lap belt and the sash belt into a unitary device.\nLap belts fasten across the lap of the occupant. These were the earliest type of seat belt, and are still standard in passenger seats of airplanes. The full force of an impact is absorbed along a relatively small area across the waist of the occupant and may cause damage to internal organs. A “sash” or shoulder harness is a strap that goes diagonally over the vehicle occupant's outboard shoulder and is buckled inboard of his or her lap. The sash belt restrains the occupant over a larger area than the lap belt. The sash belt may attach to the lap belt tongue, or it may have a tongue and buckle completely separate from those of the lap belt.\nThe 3-point belt is similar to the combination of separate lap and sash belts, but is provided as a single unit. Like the separate lap-and-sash belt, in a collision the 3-point belt spreads out the energy of the moving body over the chest, pelvis, and shoulders.\nIn most road vehicles, the three point seat belt is now standard. It comprises a tough strip of material, typically woven nylon, which is fastened to the car below the height of the seat on one side of the passenger, and is wound onto a spring loaded reel or drum that is above shoulder height or at least is threaded over a pulley that is above shoulder height and then drawn down to a drum or real that takes up spare seatbelt. Typically, a buckle that is coupled to the seatbelt is pulled down, across the passenger, from the shoulder to the other side of the passenger near the passenger's waist, where it engages a clasp provided for that purpose. Typically, pressure on a release lever on the clasp, releases the buckle, and the drum winds the belt up. In some models, the seatbelt itself is pulled down and across the occupant and is shackled by a releasable rod or loop.\nSuch seatbelts restrain the seated passenger with a strap that crosses the thighs from where a fixed end of the seatbelt is fastened to the chassis on one side, through an eye of a slider that is provided with a buckle that is slid into and held by the clasp, and then up across the body and over the shoulder to the reel or pulley wheel.\nOn impact or sudden stopping, the reel jams and prevents the seatbelt from being pulled out further, and typically is configured to retract the seatbelt. Such seat belts are reliable, convenient, and easily implemented to restrain the passenger, and easily released to allow the passenger to exit the vehicle.\nSeatbelts are mandated by law in many countries, for both front-seat and back-seat passengers of many vehicles. In some jurisdictions passengers in taxicabs and the like, are not required to wear seatbelts, but the safety advantages of so doing are nevertheless clear.\nUntil the 1980s, three-point belts were commonly available only in the front outboard seats of cars; the back seats were only often fitted with lap belts. Evidence of the potential of lap belts to cause separation of the lumbar vertebrae and the sometimes associated paralysis, or “seat belt syndrome”, led to progressive revision of passenger safety regulations in nearly all developed countries to require 3-point belts first in all outboard seating positions and eventually in all seating positions in passenger vehicles.\nMost modern seat belts are stowed on spring-loaded reels called “retractors” equipped with inertial locking mechanisms that stop the belt from extending off the reel during severe deceleration. Seatbelts in many newer vehicles are also equipped with “pretensioners” and/or “Webclamps”. Pretensioners preemptively tighten the belt to prevent the occupant from jerking forward in a crash. Some systems also pre-emptively tighten the belt in fast accelerations and strong decelerations even if no crash has happened. Webclamps clamp the webbing in the event of an accident and limit the distance the webbing can spool out (caused by the unused webbing tightening on the central drum of the mechanism). These belts also often incorporate an energy management loop (“rip stitching”) in which the lower part of the webbing is looped and stitched with a special stitching designed to “rip” at a predetermined load, which reduces the load transmitted through the belt to the occupant, reducing injuries to the occupant.\nAs with adult drivers and passengers, the advent of seat belts was accompanied by calls for their use by child occupants, including legislation requiring such use. Generally children using adult seat belts suffer significantly lower injury risk when compared to non-buckled children. Nevertheless, it will be appreciated that three point seatbelts are standard fixtures that are appropriate for restraining adults of various shapes and sizes. They are, however, engineered for adults, and are thus typically too big for small children. Adult 3 point seatbelts are generally inappropriate for children since the lower strap crosses the stomach and is consequently uncomfortable and potentially dangerous, and the diagonal shoulder strap invariably misses the shoulder of a child, and interferes with his or her neck or head. However, if the seatbelt is positioned under the arm, the seatbelt rubs against the underarm and/or neck, and is uncomfortable. Positioning the seatbelt under the arm is also dangerous, in that both sections of the belt cross the lower abdomen, and an impact may cause the seatbelt to retract across the stomach, risking rupturing internal organs.\nMany jurisdictions now advocate or require child passengers to use specially designed child restraints. Such systems include separate child-sized seats equipped with their own restraints and booster cushions or seats for use by children together with the adult restraints. In some jurisdictions children below a certain size are forbidden to travel in front car seats.\nChild safety seats (sometimes referred to as an infant safety seat, a child restraint system, a restraint car seat, or ambiguously as car seats) are seats designed specifically to protect children from injury or death during collisions. Automobile manufacturers may integrate child safety seats directly into their vehicle's design. Most commonly, these seats are purchased and installed by consumers. Many regions require children of certain ages, weights, and/or heights to use a government-approved child safety seat when riding in a vehicle. Child safety seats provide passive restraints and must be properly used to be effective. However, many child safety restraints in countries such as Canada and the United States are not used properly. To tackle this negative trend, health officials and child safety experts produce child safety videos to teach proper car seat installation to parents and caregivers. Baby car seats are legally required in many countries, including the United States, to safely transport children up to the age of 2 or more years in cars and other vehicles.\nBooster seats are recommended for children that are too big for baby seats, until they are big enough to properly use a seat belt without them. Some jurisdictions require booster seats to be used by law. Booster seats raise the child and allow the seat belt to sit firmly across the collar bone and chest, with the lap portion fitted to the hips. If the seat belt is not correctly positioned across the collar bone and the hips, it will generally ride across the neck and the stomach and cause internal injuries in the event of a collision. Booster seats also facilitate the child looking out of the window and since the child is more comfortable and has better vision, are generally accepted by children. Booster seats are required by law in the US, up to the age of 9 or a weight of 90 pounds (41 Kg).\nHowever, booster seats are large and bulky devices. Where a back seat of a vehicle is sometimes used to transport adult passengers and sometimes for transporting children, the booster has to be removed and replaced. Booster seats are space consuming and their storage in the trunk of a vehicle, minimizes the amount of groceries and other goods being transported. They are an additional expense for grandparents and other occasional drivers of children. When renting cars, booster seats may need to be rented separately, which is an additional expense, tempting parents to economize. For these and similar reasons, booster seats are often not available when required.\nVertically challenged adults may feel uncomfortable sitting on a booster, as may shorter children being transported to social events attended by taller peers.\nU.S. Pat. No. 3,136,579 to Hunter, is an early example of a booster seat that raise the child to allow him or her to look out of the windows, but is bulky and cumbersome.\nU.S. Pat. No. 3,380,775 to Offenbacher describes another bulky booster seat solution.\nU.S. Pat. No. 3,265,438 to Regan et al. describes a pneumatically inflatable seat that can be deflated for storage. Pneumatic solutions are complicated and unreliable, and this never really caught on.\nU.S. Pat. No. 3,547,489 to Grieser and U.S. Pat. No. 4,231,613 to Jonasson describe booster seats that have side loops for engaging a seat belt. The devices are fairly bulky. They are designed for the more or less obsolete lap straps, but not for a three point seatbelts.\nU.S. Pat. No. 4,291,915 to Cox describes a booster that doubles up as a suitcase and can be used for storing objects therein, going someway to overcome the inherent space requiring disadvantage of booster seats. The device described is somewhat complicated and therefore expensive to manufacture, and it is not really a practical storage solution for more than its straps.\nU.S. Pat. No. 4,461,510 to Cunningham, titled “Childs safety device for vehicle” describes a booster with a strap that goes behind the child and engages the shoulder strap of a three point seat belt. The strap may be adjusted in effective length for different sized children, but the width of the seat is fixed, and the leg straps may not be comfortable for children that are significantly narrower than the booster seat. Additionally, the booster seat is fairly cumbersome.\nU.S. Pat. No. 5,308,116 to Zawisa, describes a device that engages the shoulder strap of a three point seatbelt at an adjustable height suitable for smaller passengers. However, the system described therein is only suitable for outside seated passengers, whereas the middle seat is actually the safest for backseat passengers. Furthermore, although offering a solution for larger children, the lap engaging section of the strap will generally be too high for smaller passengers that are, nevertheless, too large for an infant seat.\nU.S. Pat. No. 5,468,020 to Scime describes an auxiliary harness that engages the seatbelt and adapts it to the child. Scime's system is different from the typical booster that raises the child, but is nevertheless, fairly bulky and thus inconvenient.\nU.S. Pat. No. 5,609,367 to Eusebi et al. describes an adjustable three point restraint system that is suitable for both children and adults. It is part of the vehicle and significantly increases manufacturing costs of the vehicle. It is not retrofittable to standard vehicles provided with conventional three point seatbelts.\nU.S. Pat. No. 5,733,004 to Celestina Krevh describes a fairly simple compact device for adjusting the shoulder strap of a three-point seatbelt and making it appropriate for use by a child. The device is designed for use with a booster seat. The device described is fairly compact, but when used with a booster seat, the space required for the booster seat makes the solution bulky and cumbersome and inappropriate for occasional use.\nU.S. Pat. No. 5,797,654 to Stroud describes a belt buckle/tether strap for use with a booster seat that combines the advantages and disadvantages of U.S. Pat. No. 4,461,510 to Cunningham and U.S. Pat. No. 5,733,004 to Celestina Krevh.\nU.S. Pat. No. 5,845,967 to Kane and U.S. Pat. No. 7,066,536 to Williams describe other bulky booster and children's seats.\nWO9700184 to Smith describes a safety belt restraint device that holds the lap belt down and has a flatter profile than the conventional booster seat. Two embodiments are described. The first one holds the belt down between the legs and over each leg separately. The second is a sort of tray that engages the seatbelt. It is appropriate for a child whose width is not much narrower than the tray. The pelvis of a smaller child will not be adequately protected.\nU.S. Pat. No. 7,703,806 to Bell et al. entitled “Seat Belt System for a Motor Vehicle” describes a system for adjusting the shoulder strap to the height of the passenger. However, the lap strap remains configured to the width of an adult and is thus inappropriate for use by smaller children as it risks damage to a narrow pelvis in impact.\nUSSN 2012/0019033 to Kelly describes an inflatable booster. By deflating the device, it may be stowed without significantly filling storage resources needed for luggage, groceries and the like. However, as described hereinabove with reference to U.S. Pat. No. 3,265,438 to Regan et al., the inflatable system is not really practical for occasional use, by grandparents and the like. Some adults are fastidious about using the mouth nozzle to inflate a booster seat, not knowing who previous users were. Inflatable seats may take two or three minutes to inflate, and require good lungs, and are thus problematic for many people such as asthmatics, the elderly and heavy smokers, for example. Additionally, the deflated seat itself takes up space.\nCZ 2007,466 describes yet another booster seat.\nDespite the plethora of developments in safety harnesses and boosters, there remains a need for a compact alternative to traditional booster seats for adapting a three point adult seat belt for use by smaller passengers. Embodiments of the present invention address this need."}
{"text": "Solid wood doors have been manufactured for many years with the same assembly principle.\nIn the last few years, the insulation of panels has become a standard step in the manufacturing of doors. The insulation of panels helps eliminate the problems of condensation on the interior of the panels, a location where the panels are the thinnest.\nIn construction, the R-value is the measurement of a material's capacity to resist heat flow from one side to the other. In simple terms, R-values measure the effectiveness of insulation and a higher number represents more effective insulation. Despite the insulation of panels, the global R value of a solid wood slab door remains relatively low.\nIn order to create a solid wood slab door with a higher R value, a system for insulating solid wood is herein provided."}
{"text": "Prostacyclin derivatives and analogs are useful pharmaceutical compounds possessing activities such as platelet aggregation inhibition, gastric secretion reduction, lesion inhibition, vasodilation, and bronchodilation.\nTreprostinil is a synthetic prostacyclin derivative currently marketed as an active pharmaceutical ingredient (API) for its ability to inhibit pulmonary arterial hypertension under the trade name Remodulin®. Treprostinil was first described in U.S. Pat. No. 4,306,075.\nProstacyclin derivatives are traditionally synthesized using a variety of methods that are described in J. Org. Chem. 2004, 69, 1890-1902, Drug of the Future, 2001, 26(4), 364-374, U.S. Pat. Nos. 4,306,075; 6,441,245; 6,528,688; 6,700,025; 6,765,117; 6,809,223 and U.S. patent application publication nos. 2009/0163738, 2011/0319641 A1, as well as Canadian patent application publication no. 2710726 A1. The entire teachings of these documents are incorporated herein by reference in their entireties. Also disclosed in these publications are processes and intermediates useful for the preparation of Treprostinil. However, the methods of these teachings suffer from one or more problems including toxic oxidation reagents, reduced yields, elevated levels of impurities, poor scalability, and numerous chromatography steps to purify intermediates and final products. Thus, there remains a need for safe, scalable, efficient, and economical processes for the preparation of Treprostinil."}
{"text": "1. Field of the Invention\nThe invention relates to an improved flow conditioner used in tubular conduits carrying single phase fluids. In particular, the invention minimizes metering errors by producing fully developed velocity profile, fully developed turbulence structure, and substantially eliminating swirl of fluids flowing in a conduit.\n2. Description of the Related Art\nThe North American natural gas industry produces, transports and distributes approximately 700 billion cubic meters of gas each year (25 trillion standard cubic feet). The Western European market transports and distributes 250 billion cubic meters of gas each year (9 trillion standard cubic feet). Because of the importance of gas measurement for industry operations and fiscal accountability, it is essential that metering be accurate, reliable, and cost efficient over a range of conditions.\nAll of this gas is measured at least once, and most of it several times, in meter sizes ranging from 25-900 mm (1-36 inches), at pressures from below atmospheric to 14 MPa (2,000 psi) , at temperatures from 0.degree.-100.degree. C. (32.degree.-212.degree. F.) , with several types of meters. Large volume metering stations utilize either concentric, square-edged, flange-tapped orifice meters or gas turbine meters.\nFor over sixty years, the concentric orifice meter has remained the predominant meter of choice for natural gas production, large volume gas flow and chemical metering applications. In fact, it is currently estimated that over 600,000 orifice meters are being used for fiscal measurement applications associated with the petroleum, chemical and gas industries in North America.\nAll flowmeters are subject to the effects of velocity profile, swirl and turbulence structure of the flowing fluid being measured. Meter calibration factors or empirical discharge coefficients are valid only if geometric and dynamic similarity exists between the metering and calibration conditions or between the metering and empirical data base conditions (i.e., fully developed flow conditions exist). In fluid mechanics, this is commonly referred to as the Law of Similarity.\nThe classical definition for fully developed turbulent flow is stated by Hinze as follows:\nFor the fully developed turbulent flow in the pipe the mean-flow conditions are independent of the axial coordinate, x and axisymmetric, assuming a uniform wall condition. PA1 a swirl-free, axisymmetric flow with time average velocity profile and turbulence structure having values approximating those found in fully developed flow and independent of the axial coordinate.\nFrom a practical standpoint, fully developed flow implies the existence of a swirl-free, axisymmetric time average velocity profile in accordance with the Power Law or Law of the Wall prediction. However, fully developed turbulent flow requires equilibrium of forces to maintain the random \"cyclic\" motions of turbulent flow. This in turn requires that the velocity profile, turbulence intensity, turbulent shear stress, Reynolds stresses, etc., are constant with respect to the axial position.\nUnfortunately, fully developed pipe flow is only achievable after considerable effort in a research laboratory. To bridge the gap between research and industrial applications, reference is made to the term pseudo-fully developed flow defined as:\nStated another way, Pseudo-fully developed flow exists when the slope of the orifice meter's discharge coefficient deviation asymptotically approaches zero as the axial distance from the orifice meter to the upstream flow conditioner increases. Of course, this assumes that the empirical discharge coefficient baseline was conducted under fully developed flow conditions.\nIn the industrial environment, multiple piping configurations are often assembled in series generating complex problems for organizations that write standards and flow metering engineers. The challenge is to minimize the difference between actual or \"real\" flow conditions in a pipeline and the vertical or research-achievable \"fully developed\" flow conditions, on a selected metering device's performance to minimize error. One of the standard error minimization methods is to install a flow conditioner in combination with upstream straight lengths of pipe to \"isolate\" the meter from upstream piping disturbances. Present domestic and international measurement standards provide specifications for upstream straight pipe lengths and flow conditioners upstream of orifice meters. See, e.g., American National Standard Institute (ANSI) (ANSI 2530) and International Standards Organization (ISO) (ISO 5167). Unfortunately, there is considerable disagreement over straight length requirements between ANSI and ISO.\nWith respect to installation effects and the near term flow field, the correlating parameters which affect similarity vary with meter type and design. However, it is generally accepted that a concentric, square-edged, flange-tapped orifice meter exhibits a high sensitivity to time average velocity profile, turbulence structure, bulk swirl and tap location.\nIn North America, current design practices utilize short upstream piping lengths with a specific flow conditioner, American Gas Association (A.G.A.) tube bundles, to provide \"pseudo-fully developed\" flow in accordance with the applicable measurement standard (ANSI 2530/A.G.A. Report No. 3/API (American Petroleum Institute) MPMS Chapter 14.3). Most North American installations consist of 90 degree elbows or complex header configurations upstream of the orifice meter. Tube bundles in combination with piping lengths of seventeen pipe diameters (17,D) have been installed in an effort to eliminate both swirl and distorted velocity profiles. Ten diameters (10*D) of straight pipe is required between the upstream piping fitting and the exit of the tube bundle, and seven diameters (7,D) of straight pipe is required between the exit of the tube bundle and the orifice meter.\nIn Western Europe, two design practices are currently employed to provide \"pseudo-fully developed\" flow in accordance with the applicable measurement standard (ISO 5167)--long upstream piping lengths with or without flow conditioners. Most Western European installations consist of complex header configurations upstream of the orifice meter. Piping lengths of one hundred pipe diameters (100*D) without flow conditioners or piping lengths of forty-two pipe diameters (42*D) in combination with flow conditioners have been installed in an effort to eliminate both swirl and distorted velocity profiles.\nThree types of flow conditioners have been generally utilized in Western Europe--tube bundles, Zanker and Sprenkle designs. Twenty diameters (20*D) of straight pipe is required between the upstream piping fitting and the flow conditioner, and twenty-two diameters (22*D) of straight pipe is required between the flow conditioner and the orifice meter.\nThe optimal flow conditioner should achieve a range of design objectives including: a minimal deviation of empirical discharge coefficient (or meter calibration factor) for both long and short pipe lengths; low permanent pressure loss across the flow conditioner (i.e., low \"head ratio\"); a low fouling rate or a low sensitivity to accumulation of foulants; elimination of swirl; and flexibility for use in both short and long straight lengths of pipe. The latter objective can be achieved by a flow conditioner that produces an axisymmetric, pseudo-fully developed time average velocity profile and turbulence structure. Additionally, it is also desirable that the flow conditioner should be subject to rigorous mechanical design and have a moderate cost of construction.\nIn the specification and claims, when the swirl angle is less than 2.degree. as conventionally measured by using pitot tube devices, swirl is regarded as substantially eliminated. Further, when the empirical discharge coefficient or meter calibration deviation for both short and long piping lengths is about 1/10 of 1% it is assumed to be at a \"minimum\".\nThe ISO and A.G.A. designs, shown in FIGS. 1A and 1B respectively, are intended to eliminate swirl. Both designs include a bundle of tubes having the same length and diameter. For the A.G.A. design (Figure 1B), the length of the bundle must be at least ten times the tube diameter. For meter runs larger than 75 mm (3 inches) the bundle typically consists of nineteen tubes arranged in a circular pattern with a bundle length of two to about three pipe diameters. For smaller meter runs, the bundle consists of seven tubes arranged in a circular pattern with a bundle length of three pipe diameters. For both the ISO and A.G.A. designs, permanent pressure loss is low, mechanical design is rigorous, cost of construction is low, fouling rate is low, and swirl is eliminated. However, the performance of these devices for minimal deviation from the empirical discharge coefficient for both short and long piping lengths is unacceptable. Also, velocity profile and turbulence structure measurements have shown that both A.G.A. and ISO designs cannot produce pseudo-fully developed flow conditions within reasonable distances due to their high porosity and constant radial resistance. This is shown by the instability in the coefficient performance graphs, FIGS. 19 and 21.\nThe Sens & Teule flow condition as shown in FIG. 2 is designed to isolate piping disturbances from flow meters. The design consists of a bundle of tubes of different lengths and diameters arranged in a circular array. Permanent pressure loss is high, cost of construction is high, and prototype designs are rigorous and complex. While swirl is eliminated, the fouling rate of this design is unknown. It has been reported that the device exhibits pseudo-fully developed time average velocity profile and turbulence structure for short piping lengths. Geometric scaling of the device is a problem, when considering a range of pipe sizes.\nFIGS. 3 and 4, respectively, show the Etoile and Air Moving and Conditioning Association (AMCA) vane-type of swirl eliminator. The Etoile design consists of three flat plates of equal length and width assembled in a star-shaped pattern around a central hub. While these designs eliminate swirl, it is known that the Etoile design does not produce pseudo-fully developed flow conditions in reasonable distances. Similarly, the AMCA design (FIG. 4) was not intended to produce a pseudo-fully developed flow.\nFIG. 5 shows an example of screens or wire gauze assembled in an egg-crate fashion within a pipeline. Fine mesh screens are impractical in an industrial environment due to high permanent pressure loss, non-rigorous mechanical construction, and high fouling rates.\nPerforated plates, such as the Sprenkle design, shown in FIG. 6, were designed to isolate piping disturbances from flow meters for measuring steam flow. The design consists of three perforated plates spaced one diameter apart and connected by rods. Each plate has a porosity of about fifty percent with regularly distributed perforations in a specified hexagonal pattern. The size of the perforations is about five percent of the pipe diameter. While the designs eliminate swirl, cost of construction is high, design is rigorous and complex, permanent pressure loss is very high, and fouling rate is moderate. Further, performance for minimal coefficient deviation for short piping lengths is unacceptable. Finally, the design is thought to not produce pseudo-fully developed flow conditions for short piping lengths due to its almost constant radial resistance.\nThe Bellinga design, shown in FIG. 7, is a modified Sprenkle design that suffers much the same shortcomings as Sprenkle.\nThe Zanker design, shown in FIG. 8A, was designed to isolate piping disturbances for the purpose of pump efficiency testing. The device consists of a perforated plate connected to a downstream grid or egg crate construction. The plate includes 32 holes of five different diameters, each hole having a specified location. Permanent pressure loss for this device is high as is cost of construction. While the design eliminates swirl, the design does not provide minimum deviation from empirical discharge coefficient for both short and long piping lengths. Therefore, the design is thought to not produce pseudo-fully developed flow conditions for all piping configurations.\nThe Akashi design, sometimes referred to as the Mitsubishi design, as shown in FIG. 9, consists of a single perforated plate with 35 holes. The hole size is 13 percent of the pipe diameter and the perforated plate thickness is equal to the hole diameter. The plate has a porosity of approximately 59 percent. Hole distribution is dense toward the (center of the pipe) and sparse around the periphery (pipe wall). The upstream inlets of the holes are beveled. While the device produces a low permanent pressure loss and mechanical design is rigorous and simple, performance for minimal deviation from empirical discharge coefficient for both and short and long piping lengths are unacceptable since the design calls for almost constant radial resistance. The design is further thought to not produce pseudo-fully developed flow conditions for all piping configurations because it does not provide minimal deviation from the empirical discharge coefficient for both short and long pipe lengths.\nThe Laws device, shown in FIG. 10, is also a single perforated plate, but with 21 holes. The plate thickness is approximately 12 percent of the pipe diameter (D) and the plate has a porosity of about 51 percent. The holes are arranged in circular spaced arrays around a central hole. The first and second arrays have 7 and 13 holes respectively. Hole size is largest in the middle of the pipe, 0.1924*D, and decreases in size to the first circular array, 0.1693*D, and further in size to the second array, 0.1462*D. The pitch circle diameter for the first and second array are about 46 and 84 percent of the pipe diameter respectively. Upstream inlets of the holes may be beveled. Once again, the performance for minimal deviation from the empirical discharge coefficient for short piping lengths is unacceptable, but acceptable for long piping lengths. The design can produce axisymmetric, pseudo-fully developed conditions only for long piping lengths.\nA further development of a device shown in U.S. Pat. No. 5,255,716 to Wilcox is the K-Labs Mark V. The patent shows a flow conditioner comprising tubular passages with the area between specific tubes blocked. While permanent pressure loss is low and mechanical design is rigorous and simple, and swirl is eliminated from most piping configurations, the performance for minimal deviation from the empirical discharge coefficient for short piping lengths is unacceptable. Therefore, the design is thought to not produce pseudo-fully developed flow conditions for all piping configurations.\nWhat is yet required is a flow conditioner for use with flow meters to provide measurements that are sufficiently accurate for industrial and fiscal applications. The flow conditioner should achieve all the design criteria stated above including elimination of swirl, and achievement of a minimal deviation from the empirical discharge coefficient or meter calibration factor for both short and long straight lengths of pipes by the production of a pseudo-fully developed time average velocity profile and turbulence structure. Further, the device should have a low permanent pressure loss (head ratio) across the flow conditioner, low fouling rate or insensitivity to foulant accumulation. Finally, the device should be subject to a rigorous mechanical design and should have a relatively moderate cost of construction."}
{"text": "This invention relates to a display hanger and, more particularly, to a hanger used to suspend merchandise from a horizontal shelf for viewing by store customers. Still more particularly, the hanger is adapted to hold plastic strip merchandisers of the general type disclosed in Fast U.S. Pat. No. 4,546,943.\nMost conventional methods of supporting plastic strip merchandisers orient the strip, merchandise attached thereto, and a corresponding label such that they face into the customer aisle provided between columns of shelves. As a result, the visibility of the merchandise is diminished because the flow of customer traffic is perpendicular to the direction in which the strip merchandisers are facing."}
{"text": "1. Field of the Invention\nThe present invention relates to a tilt detection apparatus, a hologram apparatus, a tilt correction method for a medium, and a tilt correction method for a hologram medium.\n2. Description of the Related Art\nAmong hologram media adapted to record digital data as holograms is a photosensitive resin (e.g., photopolymer) sealed between glass substrates.\nTo record digital data on a hologram medium as a hologram, a laser beam from a laser device is first split into two beams by a PBS (Polarization Beam Splitter). Then, one of the two beams (hereinafter referred to as reference beam) and a laser beam (hereinafter referred to as data beam) produced by the other beam irradiating an SLM (Spatial Light Modulator) having digital data in the form of a two-dimensional contrast image pattern, which beam reflects information of the two-dimensional contrast image pattern, are irradiated onto a hologram medium at a given angle, and thereby the digital data is recorded onto the hologram medium.\nMore specifically, the photosensitive resin making up the hologram medium has a finite number of monomers. When the laser beam (hereinafter referred to as laser beam) made up of the reference and data beams is irradiated thereonto, the monomers change into polymers correspondingly with the energy determined by the light intensity of the laser beam and the irradiation time. As a result of the transformation of the monomers into polymers, interference fringes, made up of polymers, are formed correspondingly with the laser beam energy. Therefore, as a result of the formation of such interference fringes in the hologram medium, digital data is recorded as a hologram. Later, remaining monomers migrate (diffuse) to those locations that have had monomers consumed. Further, by irradiating the laser beams again, changes of such monomers into polymers are iterated. FIG. 2 schematically illustrates how monomers transform into polymers in response to the laser beam energy in the hologram medium.\nWhen there is a large amount of digital data to be recorded in the hologram medium, so-called angle-multiplexed recording may be performed which forms a number of holograms by varying the incidence angle of the reference beam onto the hologram medium. For example, a hologram formed on the hologram medium is called a page, whereas a multiplexed hologram made up of a number of pages is called a book. FIG. 3 schematically illustrates a book and pages in the angle-multiplexed recording. As shown in FIG. 3, for a single book in the angle-multiplexed recording, e.g., ten pages of holograms are formed by varying the incidence angle of the reference beam. Thus, the recording of digital data onto a hologram medium by the angle-multiplexed recording allows for the recording of a large amount of digital data.\nMeanwhile, to reproduce digital data from the hologram medium, the reference beam is irradiated onto the interference fringes representing the digital data at the same incidence angle as when the interference fringes were formed, and the reference beam diffracted by the interference fringe (hereinafter referred to as reproduction beam) is received by an image sensor or other means. The reproduction beam received by the image sensor or other means produces a two-dimensional contrast image pattern representing the above-mentioned digital data. Then, the digital data can be reproduced by demodulating this two-dimensional contrast image pattern by a decoder or other means.\nAs such, when digital data is reproduced from a hologram medium, a two-dimensional contrast image pattern is reproduced from the reproduction beam. Hence, the reproduction beam at an image sensor or other means must have a light intensity equal to or above a given level to allow the reproduction of the two-dimensional contrast image pattern. To have the reproduction beam have at least the given level of light intensity, therefore, the interference fringes diffracting the reference beam must have a specified or higher value of diffraction efficiency, which indicates the ratio of the reproduction beam light intensity to that of the incident reference beam. It is to be noted that the specified value of diffraction efficiency is such a value that the reproduction beam has the given level of light intensity.\nFIG. 4 is a diagram showing an incidence angle θ between the reference beam (hereinafter called a recording reference beam) and the data beam when forming the interference fringes in a hologram medium. FIG. 5 is a diagram showing an angle θ between a reference beam (hereinafter called a reproducing reference beam) and the reproduction beam (hereinafter called a reproducing angle θ) when reproducing digital data from the interference fringes of FIG. 4.\nWhen reproducing digital data from the hologram medium, letting Δθ, n, t, and λ be the angle difference between the incidence angle θ and the reproducing angle θ, the refraction index of the hologram medium, the thickness of the hologram medium, and the wavelength of the laser beam emitted from a laser device respectively, Δθ is expressed as (λ×√(n2−sin2θ))/t×sin θ×cos θ. If n=1 and sin θ=1, then Δθ=λ/t, which is found to be no greater than the wavelength of the laser beam. If the laser beam is, e.g., a helium neon laser, its wavelength λ is 633 nm, and Δθ is very small. As a result, it is understood that the reproducing angle θ for reproducing digital data from the hologram medium needs to be within such an angle difference Δθ of the incidence angle θ for when the interference fringes representing the digital data are formed. That is, the reproducing reference beam needs to be incident on the hologram medium at a very accurate reproducing angle θ. And the hologram medium needs to be placed in such a position (hereinafter called a reference position) that the reproducing reference beam is made incident thereon at the very accurate reproducing angle θ.\nMeanwhile, it is required for the sake of convenience, flexibility, and the like that the hologram medium is configured to be attachable to and detachable from a hologram apparatus able to record or reproduce holograms. For example, when a hologram medium on which a hologram apparatus recorded digital data as a hologram is required to be played back by another hologram apparatus, if the hologram medium is attachable, the requirement can be easily fulfilled (refer to, e.g., Japanese Patent Application Laid-Open Publication Nos. 2004-177958 and 2004-272268).\nHowever, in the case of an attachable hologram medium, the hologram medium may be tilted with respect to the above reference position depending on the mechanical configuration and the like of a body on which the hologram medium is mounted. The tilt of the hologram medium may cause the angle difference between the incidence angle θ and the reproducing angle θ to be greater than the angle difference Δθ thereby making the reproducing reference beam not be accurately incident and disabling the reproduction of digital data. Or, the reproducing reference beam is not accurately incident on a desired hologram but on another hologram, and false digital data may be reproduced. Also, when recording digital data onto the hologram medium, the digital data may not be recorded as a hologram at a desired position in the hologram medium.\nIn particular, where the hologram medium is shaped like a disk, since being thin, it is difficult to hold the hologram medium horizontal with respect to the reference position, and the hologram medium may be tilted. Moreover, the disk-shaped hologram medium, in order to make it attachable to a hologram apparatus, may be mechanically processed to have, e.g., a central opening therein. And the hologram apparatus is provided with a support that fits detachably the central opening and with a table to rotate hologram media to form a number of holograms. As a result, due to distortion that occurs in mounting hologram media on the hologram apparatus or centrifugal force exerted on the hologram media when rotating, the hologram media may be tilted with respect to the reference position."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method of forming a fusion welded butt joint between thermosplastic pipe sections.\n2. Description of the Prior Art\nPipelines made of medium to high density polyethylene or similar thermoplastic material are commonly used for transporting natural gas or the like. While, upon original installation a joint may be made by applying heat direct to the ends of pipe sections to fuse them directly together to form a butt joint, in certain applications, as for instance in repairing pipe previously buried underground, the pipe is so restricted against movement that a direct butt joint cannot be made. This is because, when a section of replacement pipe is installed in close fitting relationship within a gap formed between the ends of the pipe left when a damaged section of pipe has been removed, there is insufficient space between the confronting pipe ends to receive a heating rod for melting the pipe to its molten, or fusing, temperature thereby achieving direct fusion. In these instances there have been efforts to joint the pipe sections together by a thermoplastic coupling collar which receives the pipe sections within its opposite ends. The collar incorporates a helical heating coil formed by continuous wraps of coil spaced longitudinally along the inner surface thereof. Energization of the coil raises the temperature of the interface between the coupling and the abutting pipe sections to its molten level to cause fusion. This type of fusion welded joint is characterized by a number of significant shortcomings.\nMore particularly, when the interface between the coupling collar and the pipe sections is heated to a molten state, the wraps or lops of the heating coil tend to float in the molten material frequently resulting in some becoming disposed in more closely spaced relationship than others. In the areas of the more closely spaced wraps there will be a concentration of heat while in the area of the wraps of greater spacing there will be an absence of heat resulting in an uneven temperature profile along the length of the heating coil. Oftentimes the location of the elevated temperature will be at the butt joint interface resulting in the pipe wall becoming more molten than the remainder of the length of collar. This results in the thermoplastic material forming the wall of the pipe adjacent the butt joint losing its structural rigidity causing it to deform or flow. As a consequence, the end portions of the pipe wall adjacent the butt joint will sag radially inwardly, and the tendency of the wraps, or loops, of the helical heating element to migrate toward the sagging area will be enhanced. The resulting concentration of heat at the butt joint compounds the sagging problem, accelerating the flow of molten material and the migration of the heating coil loops toward the butt joint area, all before the interface throughout the remainder of the collar has reached its fusion temperature. Since the heating coil must remain heated until fusion has been achieved between the collar and pipe sections for some distance from the pipe ends, the temperature in the sagging areas will continue to elevate creating \"hot spots\" where the thermoplastic pipe walls will become burned thus leading to deterioration of the structure of the plastic itself. These shortcomings are compounded by the fact that the pipe ends are frequently of an irregular cut thus leaving gaps of as much as 1/8 inch one one side or the other when the ends are abutted together for fusion. The amount of parent material required from the wall of the pipe to fill such gaps during fusion is limited thus frequently leading to the wall of the joint in that area being relatively thin thereby resulting in the section being heated to an even higher temperature during fusion. The fusion welded joint which develops under this method is imperfectly formed and susceptible to structural failure and gas leakage.\nAttempts have been made in the prior art to solve the development of such \"hot spot\" heating adjacent the butt joint by employing two separate helical coils spaced apart on opposite sides of the butt joint to thus avoid applying heat directly to the interface of the pipe ends. This is a relatively expensive expedient which is not completely effective and leaves the unheated portion at the interface of the pipe ends unfushed and susceptible to subsequent entry of fluids under pressure between the pipe wall and sleeve, a condition which in use may create stress risers thus further deteriorating the joint.\nAnother prior art attempt to preserve the structural integrity of thermally fused pipe section end portions is disclosed in U.S. Pat. No. 3,788,928, issued Jan. 29, 1974 for a \"Method of Forming a Lap Joint Between Tubular Articles of Thermoplastic Material\". As its title indicates, the patented method is directed to a lap joint. It utilizes a reinforcing sleeve of a material having a melting temperature higher than that of the thermoplastic pipe. The sleeve is inserted into one pipe section end portion to support the joint area such that when the pipe is heated only to its fusion temperature, the insert will not melt and will afford support against radially inward sagging during the fusion welding step. However, since the method is limited to an insert of dissimilar material to lap joints which requires axial movement of the pipe sections to make the joint, it is not practical for use in pipe repair as described above. Furthermore, the method itself contemplates that the temperature will not reach the fusion temperature of the sleeves, thus having the disadvantage that material of the flange thereof will not be available to fill any spacing that may exist at any location around the circumference of the pipe ends.\nOther efforts for a satisfactory joint include the construction of a piston style dynamic coupling as shown in U.S. Pat. No. 4,386,796 assigned to the assignee of this application. That coupling, which satisfactory for its intended purpose, is relatively expensive to manufacture and is not satisfactory for all applications."}
{"text": "The heat exchanger for vehicles disclosed in Japanese Unexamined Utility Model Publication No. H1-61582 is achieved by forming a heat exchanger for engine cooling water, a heat exchanger for air conditioning and other heat exchangers as an integrated unit, with each heat exchanger provided with a core constituted of a plurality of tubes and fins secured in contact with the tubes and a tube plate that covers the ends of tubes belonging to two cores collectively. In addition, a groove is formed at the circumferential edge of the tube plate, and the bottom portion of the tank main body constituted of a synthetic resin is fitted and fastened through calking at the groove.\nThe radiator illustrated in FIG. 23(a) is a so-called down-flow radiator and assumes a structure similar to that described above. In more specific terms, this radiator 100 is provided with tank main bodies 102 and 103 constituted of a synthetic resin and disposed at the top and the bottom of a core main body 101 constituted of tubes 104 and fms 105 both constituted of aluminum alloy. As shown in FIG. 23(b), the tank main bodies 102 and 103 each have a flange portion 108 which is fitted via an o-ring at a groove 107 formed at the periphery of an end plate 106 to which ends of the tubes 104 are mounted and the tank main bodies 102 and 103 are each further fastened by using calking tabs 109 formed over specific intervals at the circumferential edge of the end plate 106.\nIt is to be noted that in FIG. 23(a) illustrating the radiator 100, reference number 110 indicates an intake pipe through which engine cooling water is guided into the upper tank main body 102 and reference number 111 indicates an outlet pipe through which the engine cooling water is discharged from the lower tank main body 103. In addition, a cooling water induction port 116, which is closed off by a cap 112 having a pressure valve, for instance, is provided at the upper tank main body 102. Inside the lower tank main body 103, an oil cooler is provided, and reference numbers 114 and 115 indicate intake/outlet pipes of the oil colors.\nHowever, in the structure of the prior art described above, in which the tubes and the fins constituting the core are formed from aluminum alloy and the tank main bodies are formed from a synthetic resin, there is a problem in that they cannot be formed together. There is another problem in that the recyclability of the radiator itself is poor.\nAs a solution, a method achieved by forming the members constituting the tank portions with aluminum alloy and then the aluminum alloy tank portions are brazed together with the core in a furnace to achieve an integrated unit may be proposed. However, a problem occurs during the repair process implemented after the brazing process to repair any defective brazing occurring between the individual members constituting the tank portions by means such as torch brazing or the like that is, the brazed areas between the individual members are close to the tubes and fins, the tubes and fins become melted during the repair process.\nIn addition, while the oil cooler for cooling the automatic transmission oil (hereafter referred to as the A/T oil cooler) is mounted at the same time inside the outlet-side (lower) tank main body 103 in the radiator, if U-shaped tank plates are used, the intake/outlet pipes of the A/T oil cooler become a hindrance to the assembly work. Furthermore, while the intake/outlet pipes of the A/T oil cooler may be enclosed and brazed between the tank plates, this method poses problems in that the shapes in the vicinity of the insertion holes for the intake/outlet pipes are bound to become complicated and in that good brazing is not achieved for the intake/outlet pipes, the tank plates and the like.\nAn object of the present invention is to provide a heat exchanger with a structure that allows integrated brazing, that achieves an improvement in the assemblability in the mounting of the A/T oil cooler and also achieves good overall assemblability and good recyclability."}
{"text": "1. Field of the Invention\nThis invention relates to a voltage detecting device for detecting a voltage provided at a predetermined part of an object under measurement, and more particularly to a voltage detecting device utilizing the principle that the polarization of light passing through electro-optical material is changed by a voltage provided to that material\n2. Background Information\nA voltage detecting device for detecting, in a time resolution of sub-pico seconds, the voltage of an ultra-high-speed photo detector, semiconductor switch or high-speed electronic device is well known in the art.\nIn FIG. 6(a) is shown a perspective view of a voltage detecting device of the above-described type, which was disclosed by Janis A. Valdmanis et al, \"IEEE Journal of Quantum Electronics,\" Vol. QE-19, No. 4, pp. 664-667 (published in Apr. 1983). Shown in FIG. 6(b) and FIG. 6(c) are a top view and a front view of the voltage detecting device shown in FIG. 6(a).\nIn the voltage detecting device as shown in FIG. 6(a), an electro-optical material 50 of lithium tantalate (LiTaO.sub.3) is cut perpendicular to the C-axis, an aluminum strip line 52 is provided on a surface 51 of the electro-optical material 50 which is perpendicular to the C-axis, and a predetermined part of an object 53 under test is connected to the strip line 52.\nWhen the object 53 under test is an ultra-high-speed photo detector, for example, with a predetermined part of the photo detector connected to the strip line 52, a voltage pulse VP having a pulse width of the order of ten pico seconds outputted from photodetector 53 moves along the strip line 52 at a speed V.sub.O as shown in FIG. 6(c). As a result, an electric field E is applied to that part of the electro-optical material 50 which is just below the strip line 52 along which the voltage pulse VP moves. The refractive index of that part of material 50 is changed as a result of the electric field. Accordingly, referring to FIG. 6(b), when a linearly polarized light beam PB is applied to one side 54 of the electro-optical material 50 in such a manner that it forms an angle .theta. with the longitudinal axis A--A of the strip line 52, the speed component, V cos .theta., of the light beam PB taken along the axis A--A is equal to the speed V.sub.O of the voltage pulse VP. Thus, the light beam PB passes through the electro-optical material 50 following the refractive index change caused by the voltage pulse VP. The polarization of the light beam will be changed by the well known Pockels effect as it passes through material 50. The light beam will be outputted, as a transmitted light beam, from the opposite side 55 of the electro-optical material 50. If the change in polarization of the transmitted light beam can be detected, then the value of the voltage pulse VP moving along the strip 52 can be achieved without affecting the object under test.\nIn the voltage detecting device shown in FIG. 6(a), the light beam PB is applied to the electro-optical material 50 in such a manner that its speed component along the longitudinal axis of the strip line 52 is equal to the speed V.sub.O of the voltage pulse VP; that is, the light beam PB and the voltage pulse are allowed to interact with each other, so that the voltage is detected from the change in polarization of the light beam PB. In the voltage detecting device, as described above, the light beam PB is allowed to enter the electro-optical material 50 through one side 54 and to emerge from the opposite side 55. Therefore the optical beam PB and the voltage pulse interact with each other for the period of time T, which represents the time the light beam PB spends crossing the strip line 52. This relationship is given by the following formula: EQU T=W/(V sin .theta.) (1)\nwhere W is the width of the strip line 52, and V is the speed of the light beam PB in the electro-optical material 50.\nIf the voltage pulse VP is of the order of several kilovolts (kV), the polarization of light beam PB will change to the extent that it can be detected, even if the period of time T is short. However, if the voltage pulse VP is lower than several kilo-volts, then it is necessary to increase the period of time T in order to detect the polarization of the light beam PB. Since the period of time T depends on the width W of the strip line 52, as is apparent from the above equation (1), the ability to design the voltage detecting device shown in FIGS. 6(a) and (b) to increase the period of time T, and to change the polarization of the light beam PB as required, is limited to the practical width W of strip line 52.\nAccordingly, an object of this invention is to provide a voltage detecting device which can change the polarization of a light beam in an electro-optical material to the extent that it can be detected, to detect the voltage value of the voltage pulse with high accuracy."}
{"text": "Maintaining the shape and overall dimensions of the tire bead is critical for assembly and manufacture of high quality tires, however, conventional tire beads typically distort or deform during tire assembly and curing. Conventional bead insulation compositions typically have good bead wire adhesion and processability but do not maintain their shape and overall dimensions (i.e., geometry) in an uncured or green state. “Processability” is a term understood by those skilled in the art, and is typically indicated by scorch resistance, viscosity and the ability of the rubber composition to be used in bead forming equipment known to those skilled in the art.\nNumerous solutions have been adopted by those skilled in the art in an attempt to overcome the problem of distortion and deformation of the shape and overall dimensions of conventional tire beads during tire manufacturing, yet each has associated disadvantages. One solution has been to precure the green tire beads so as to provide a tire bead which will maintain its bead geometry during tire manufacturing. Conventional tire beads are typically precured in an autoclave at 285° F. for 11 minutes. A disadvantage of precuring is that it raises manufacturing costs by adding an additional manufacturing step, raising energy costs and lowering tire bead manufacturing yields.\nTwo additional solutions, which also involve preprocessing steps, have been to wrap or staple tire beads. Wrapping the beads requires significant manpower and time, and, accordingly, significantly raises manufacturing costs. The disadvantages of stapling the bead are that special automatic equipment is required involving significant capital, increased operational complexity, lowered process yields, and, accordingly, also increased manufacturing times.\nTire manufacturing involves forcing two tire beads over the ends of a tire carcass. The tire carcass and conventional tire bead is then assembled with other tire components, such as the tread and side wall components, and then vulcanized. The tire bead must have dimensional integrity such that the tire bead can be placed over the tire carcass, and turned thereover, without distorting or deforming the tire bead.\nU.S. Pat. No. 4,575,534 to Oshima et al., Pat. No. 4,540,744 to Oshima et al., Pat. No. 5,283,289 to Yamamoto et al., Pat. No. 5,098,941 to Tsuriya, and Pat. No. 3,715,266 to Winters et al. disclose conventional tire rubber compositions which contain uncrosslinked styrene butadiene rubber. In particular, the '744 patent to Oshima discloses a tire rubber composition containing a butadiene polymer with a branched polymer through a tin-carbon bond and a styrene butadiene rubber. The '289 patent to Yamamoto et al. discloses a fiber reinforced styrene butadiene containing rubber composition. The '941 patent to Tsuriya discloses a phenolic resin reinforced bead filler composition having a low amount of hexamine which reduces the amount of deterioration of the tire carcass. The '266 patent to Winters et al. discloses adding an adhesion component to conventional uncrosslinked styrene butadiene to improve the adhesion between the bead wire and the bead insulation composition. A disadvantage of these conventional tire rubber compositions is that they typically must be precured, cloth wrapped or stapled to maintain dimensional stability when used in making tire beads.\nWhat is desired, therefore, is a bead insulation composition and tire bead that is dimensionally stable and that does not have to be preprocessed, but which provides sufficient bead wire adhesion and processability. Also desired is a method of manufacturing a tire using a dimensionally stable tire bead which exhibits sufficient bead wire adhesion and processability but which is not preprocessed, and a method of making a tire bead using a dimensionally stable bead insulation composition."}
{"text": "According to the related art, various application services may be enabled by inserting additional data into an image file.\nAs example of the related art, the following two patents are discussed below: U.S. Pat. No. 6,993,196 “Digital image storage method” (hereinafter referred to as related art 1) and U.S. Pat. No. 8,340,440 “Apparatus and method for inserting additional data into image file in electronic device” (hereinafter referred to as related art 2).\nIn the specification, for convenience of explanation, an image file having additional data inserted therein is referred to as a compound file.\nIn such a way, once the compound file is created, various application services may be enabled when the image file is played. For example, when the additional data is sound source data, a sound collected while the corresponding image file is captured or a sound harmonizing with the image file may be provided while the image file is played. Alternatively, various pieces of information or information obtaining sites (for example, a Uniform Resource Locator (URL) link) relative to the image file may be provided.\nOn the other hand, in order to widely use a compound file, it becomes an important issue to rapidly determine whether a specific image file is a file having image data (e.g., an image file, a compound file having additional data, and the like) inserted therein. Furthermore, even when it is determined as a compound file, it may become an important issue to rapidly access additional data by a data processing device which processes the image file. In addition, wide usability of the compound file or compatibility between applications can be ensured by allowing the additional data to be easily inserted or removed.\nHowever, the related art 1 discloses a method for confirming whether an image file is a compound file. For this, an End of Image (EOI) is compared with an End of File (EOF), and, when the EOF is larger than the EOI, (i.e., some data is further included behind image data), the image file is determined as the compound file. However, although the data processing device handling the image file can rapidly access the EOF, it takes a lot of time and work to determine the EOI. For example, according to a type of image file, there is a case where information for a location of the EOI exists in a header, but there is also a case where information for the location of the EOI does not exist in the header. In addition, although the header has the EOI information, a process is needed to search at least the header of the image file for information on the EOI. Accordingly, a lot of time and accesses are needed for determining whether the image file is a compound file, and accordingly a lot of time is to be taken to access the additional data by the data processing device.\nMoreover, the related art 2 discloses that information for presence or absence of additional data is recorded in a header. This method is illustrated in FIG. 1.\nFIG. 1 illustrates a schematic structure of a compound file according to the related art.\nReferring to FIG. 1, a compound file 10 includes an image file and additional data 13 which is inserted into the image file. The image file includes a header 11 and image data 12. As illustrated in FIG. 1, the header 11 has important information recorded therein for handling the image file, such as a type of the image file, and information related to encoding and decoding. In addition, in the related art 2, the header 11 has information for presence of the additional data or a location of the additional data recorded therein. Accordingly, even in the related art 2, a long time may be taken to sequentially search the header 11 in order to confirm whether the image file is a compound file, that is, a longer time to be taken as a size of the header 11 becomes larger. Furthermore, when the presence of the additional data or an item related to the additional data is recorded in the header 11, there is a risk for losing information for the additional data when an event (e.g., re-encoding, resizing, changing of the image file, and the like) by which header information may be changed occurs, while the image file is handled by a data processing device. Moreover, in order to prevent this type of risk, the information on the additional data needs to be separately managed in occurrence of the event, which causes problems in wide usability or compatibility for various applications which handle corresponding compound files. In addition, there may also be a limitation of changing the header 11 when the additional data is further inserted or removed.\nTherefore, a need exists for a compound file format, a compound file creating method, and a data storage device capable of rapidly confirming whether an image file is a compound file, rapidly accessing the additional data, and easily performing insertion or removal of the additional data.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "1. Field of the Invention\nThe present invention relates to fluid filtration, and particularly to a filtration system including a vessel containing a media through which a fluid is filtered. More particularly, the present invention relates to a distributor/collector system positioned within the vessel for directing the flow of fluid through the media bed.\n2. Description of the Related Art\nConventional distributor/collector systems are positioned in a vessel with the collector spaced-apart from the distributor and positioned downstream of the distributor. A typical distributor/collector system includes an upper header as part of the distributor and a lower header as part of the collector, each running generally transversely through the vessel and generally parallel to each other. Both the upper header and lower header include a series of laterals extending radially away from the header and generally in a plane perpendicular to a longitudinal axis of the vessel. In a typical distributor/collector system, a fluid (e.g. water, etc.) enters the vessel and flows down the upper header generally along a diameter of the vessel and then away from the upper header through the upper laterals. From there, the water flows through a media bed below the upper header and into the lower laterals below the media bed. The water then is collected by the lower laterals and carried back toward the center of the vessel to the lower header, which sits generally along a diameter of the vessel, like the upper header. Finally, the water flows down the lower header and exits the vessel on generally the same side of the vessel as it entered the vessel.\nIn conventional distributor/collector systems, the flow may not be uniform throughout the media bed. The time it takes water to flow through the system at certain points within the media bed may vary from the time it takes water to flow through the system at other points within the media bed. Water which flows out the upper laterals close to the upper header also flows into the lower laterals close to the lower header and, therefore, travels a shorter distance through the system than water which travels out the upper laterals at a distance away from the upper header and into the lower laterals at a point distant from the lower header. These various flow paths through a typical system create various levels of flow through the system. A distributor/collector apparatus that provides even flow throughout the system and, thus, a high distribution efficiency (the ratio between the lowest flow and highest flow that exits the distributor or enters the collector) would be welcomed by users of filtration systems utilizing such a distributor/collector."}
{"text": "After bituminous paving materials have been removed from a roadbed, they are hereafter referred to as recycled asphalt pavement or RAP.\nIt is generally known that the majority of existing roadways, both concrete and bituminous asphalt, undergo constant repair and surface overlay with new hot mix asphalt to achieve and maintain safe and comfortable high speed riding surfaces.\nIn recent years, new equipment has been introduced to the road paving industry in the form of pavement milling or grinding machines. The science of preparing an old roadbed base for new resurfacing is now commonly referred to as milling. Both state and Federal Department of Transportation (DOT) agencies throughout the country have readily accepted the science of milling.\nThe milling of old road surfaces provides a number of advantages in preparing the old roadbed for resurfacing. Milling not only ensures a new, smooth and level base for the new hot mix overlay, but at the same time lowers the road bed height to maintain bridge deck clearances and curb and gutter depths. Grinding or milling is also beneficial in removing potholes, old cracks, joint seams, and ruts along with other surface damage that would quickly reappear in a new surface overlay if not repaired. With many of the state and Federal DOT agencies now requiring the milling of road surfaces before permitting new overlay, there is an increasing inventory of asphalt millings being generated. The piles of discarded asphalt millings are becoming problems for land use, aesthetics and the environment. Attempts at reuse have proven difficult.\nNeeds exist for new recycling methods and apparatus for asphalt paving millings."}
{"text": "Laser-processing of features (e.g., through-vias, blind vias, trenches, routs, kerfs, and other features) within one or more materials of a workpiece can be processed using parameters such as laser power, pulse repetition frequency (PRF), pulse energy, pulse width, bite size, and other parameters. In many laser processing applications, the speed or efficiency with which a feature is formed, and the quality of the feature ultimately formed, can be very sensitive to such processing parameters.\nAn example application in a compound motion laser processing machine is the processing of “rout” features: laser cut lines made up of sequential linear or circular arc segments. Traditionally, such processing is accomplished by moving the process beam at a constant velocity along a desired trajectory. For a given laser power and PRF, this provides a consistent fluence and bite size at the worksurface.\nHowever, it is possible to exceed the dynamic limits (e.g., linear stage acceleration or velocity, or galvo scan field) of a compound motion system when processing routs. For certain rout velocities, for example, a rout that reverses direction can easily produce peak linear stage acceleration beyond the system's capability."}
{"text": "The invention relates to a method of transmitting or storing, over an interference affected channel, digital audio signals, wherein transmission errors are detected at the receiving end and corrected if necessary or masked, with the masking being effected in that the interfered-with signal section is muted or replaced by a signal preceding the interfered-with signal section in the same channel or by a synchronous, not interfered-with signal section of an adjacent channel. Such a method is disclosed in German Patent DE 3,638,922.C2.\nIn digital audio signal transmissions and storage, the received and read-out audio signals, if they contain bit errors, may either be not decodable at all or not decodable in part after all error detection and error correction methods have been exhausted. In that case, it is the custom to switch to a decoder for muting over a broad band in such a way that for a certain time period the entire signal is set at 0. In the case of digital audio signals transmitted by radio, this case occurs relatively frequently at the fringes of the reception area which is extremely annoying particularly in connection with mobile reception. The same applies for audio signal storage if the tape material or audio heads are worn to a degree that exceeds a tolerance value.\nTo reduce the annoying effect of missing signals in radio transmissions, German Patent DE 3,638,922.C2 discloses a mutual offset in time between the left and right stereo channels and, if there is an uncorrectable signal interference, placing the complementary stereo information transmitted at an earlier or later point in time instead of the interfered-with original information. Although the complementary stereo information is not identical with the associated, interfered-with original information, but is connected with it only by way of left-right correlation, such a substitution is acceptable in any case for a short period of time since direction and distance perception as well as the perception of spatial relationships are subject to a certain inertia in the human ear. However, if such a masking technique is employed for a longer period of time, the stereophonic impression is lost, since the masking always covers the full bandwidth of the interfered-with signal and therefore also replaces spectral signal components that did not suffer interference.\nIn contrast thereto, it is the object of the invention to provide, in a method of the above-mentioned type, a subjectively better masking of errors which permits better reconstruction of the interfered-with original audio event and, particularly also during longer use, maintains the stereophonic impression.\nThe above and other objects are accomplished according to the invention by the provision of a method of transmitting or storing, over an interference affected channel, digital audio signals that have been subjected to data reduction to result in a reduced data, digital audio signal that is present in at least one of the time domain and spectral domain, depending on source coding, the method comprising the steps of: detecting an interfered-with signal section in the reduced data, digital audio signal at a receiving end; and masking the interfered-with reduced data, digital audio signal section by one of (a) muting only interfered-with spectral values or subbands, or groups of spectral values or subbands in the signal section and (b) replacing only interfered-with spectral values or subbands, or groups of spectral values or subbands on the signal section with a signal component from the same channel or an adjacent channel including at least one of time domain sampled values, spectral domain sampled values, scale factors, and control information, depending on the source coding of the reduced data, digital audio signal.\nThe invention is based on the consideration of intentionally muting, repeating, estimating or replacing for a certain period of time only those spectral components of a complete audio signal that have actually experienced interference. The invention takes advantage of the fact that reduced data, digital audio signals (i.e., digital audio signals subjected to known data reduction processes) are present in the time and spectral domains, depending on whether they are subdivided into subbands (subband coding) or into spectral values (transformation coding). The following masking strategies can be employed:\n1. Muting\nWith this masking strategy, only those spectral components of the interfered-with signal are muted which have actually been interfered with. Depending on the source coding method employed, these may be individual subbands or spectral values or groups of subbands or spectral values of a digital audio signal that has been subdivided into n subbands or spectral values, with n being a value equal to or greater than 1.\n2. Repeating\nWith this masking strategy, components of the same channel signal are employed as substitutes for the interfered-with signal components, which hereinafter will be called xe2x80x9crepeating.xe2x80x9d Depending on the source coding method employed, the components to be repeated may be individual subbands or spectral values or groups of subbands or spectral values of a digital audio signal that has been subdivided into n subbands or spectral values, with n again being a value equal to or greater than 1. Again depending on the source coding method employed, these components may be composed of synchronous or spectral sample values as well as control informations or scale factors. The repetition may also be effected several times.\n3. Left-Right Substitution\nIn this case, synchronous components of the audio signal that are present without interference or in processed form in the adjacent channel are employed as substitutes, which will hereinafter be called xe2x80x9cleft-right substitution.xe2x80x9d The components to be replaced may be, as in the case of repeating and depending on the source coding method employed, individual subbands or spectral values or groups of subbands or spectral values of an audio signal that has been subdivided into n subbands or spectral values, with n again being a value equal to or greater than 1. Again depending on the source coding method employed, these components may be composed of time domain or spectral domain sampled values as well as of control information or scale factors.\n4. Estimating\nWith this masking strategy, not interfered-with or processed components of the same channel signal or of the signal from the adjacent channel are utilized by way of estimation (e.g. interpolation) to determine the components required to substitute interfered-with components. Due to the probability of combinations in the time domain or the spectral domain, a conclusion can be drawn from the use of non-interfered-with or processed spectrally or timely adjacent components of the same channel signal or of a signal from the adjacent channel as to the original content of interfered-with components that must be replaced. Depending on the source coding method employed, the components to be estimated may be individual subbands or spectral values or groups of subbands or spectral values of an audio signal that has been subdivided into n subbands or spectral values, with n here again being a value equal to or greater than 1. Again, depending on the source coding method employed, these components may be composed of time domain or spectral domain sampled values as well as of control informations or scale factors. Various methods for fading the signals in and/or out can be employed for the above-mentioned masking strategies. That means that the changeover from the not interfered-with component to the replaced or muted components and/or from the replaced or muted component to the not interfered-with component is made over a certain transition time and with a certain transition function so that sudden transitions are avoided.\nAll of the above-mentioned masking strategies can be combined with one another if required.\nIn connection with the described masking strategies, the invention takes advantage of certain characteristics of the human sense of hearing. In the case of muting, repeating or estimating, the auditory characteristics with respect to timely and/or simultaneous masking thresholds are utilized to the extent that these masking measures remain substantially inaudible or hidden as long as they do not exceed a certain time, spectrum and level range.\nIn the case of the left-right substitution, advantage is taken of the fact that direction and distance perception as well as the perception of spatial relationships are subject to a certain inertia. This can be utilized to the extent that short-term changes in the display of directions and distances of sound sources and of a spatial impression remain inaudible if they are shorter than a certain length of time and do not occur too frequently. If these two conditions can be met, the correct application of the described masking strategies results in a substantial improvement over conventional, known masking techniques."}
{"text": "1. Field of the Invention\nThis invention relates to a semipermeable membrane for separating ingredients of liquid mixture by selective permeation, and more specifically, it relates to a semipermeable composite membrane for the desalination of sea water, brackish water and selectively removing or recovering pollutant or useful materials from waste water causing public pollutions, for example, from dyeing process or from electrodepositing process, and thereby capable of contributing to the completeness of treatment of such waste water under a closed system.\n2. Description of the Prior Art\nAsymmetric membranes prepared from cellullose acetate are known as the semipermeable membranes utilized industrially, for example, so-called \"Loeb\" type of those membranes as disclosed in U.S. Pat. Nos. 3,133,132 and 1,133,137. However, these membranes are poor in the resistivity to hydrolysis, bacterial attack and chemicals. Particularly, if tried to improve their permeability, their pressure resistance and durability are sacrificed, and accordingly they have not yet been put into practical use for a wide range of application although employed in a restricted field.\nExtensive studies have been made for investigating a new material free from the drawbacks of the asymmetric cellulose acetate membranes mainly in U.S.A. and Japan, and those materials partially improved the foregoing drawbacks have been obtained, such as aromatic polyamide, polyamide hydrazide (U.S. Pat. No. 3,567,632), polyamide acid (Japanese Patent Publication No. 55-37282, U.S. Pat. No. 3,878,109), cross-linked polyamide acid (Japanese patent Publication No. 56-3769), polyimidazopyrrolone, polysulfonamide, polybenzimidazole, polybenzimidazolone and polyarylene oxide. However, none of them is superior to the cellulose acetate membrane with respect to the selective separation characteristic or permeability.\nOn the other hand, composite membranes comprising ultra-thin membrane which is substantially responsible to membrane performances and a microporous substrate for supporting the ultra-thin membrane have been developed as different type of semipermeable membrane from the \"Loeb\" type. In such composite membrane, it is possible to select the most optimum materials to each of the ultra-thin membrane and microporous substrate according to the purposes, and this can increase the degree of freedom for membrane fabrications. Another advantage of those composite membranes is said that they can be stored under the dry state different from the \"Loeb\" type membranes which must be stored always in a wet condition.\nThese composite membranes are classified into two types, that is, those comprising an ultrathin membrane and a microporous substrate coated with the ultrathin membrane directly and those comprising an ultrathin membrane, a microporous substrate and a gel layer intermediating between them. Example of the former are described in U.S. Pat. Nos. 3,744,642 and 3,926,798, Japanese Patent Application Laid-Open No. 55-147106 (U.S. Pat. No. 4,227,344), U.S. Pat. Nos. 4,557,949 and 4,366,062 and Japanese Patent Application Laid-Open No. 58-24303, and examples of the latter are described in Japanese Patent Application Laid-Open No. 49-133282 (U.S. Pat. No. 4,039,440), Japanese Patent Publication Nos. 55-49524 (U.S. Pat. No. 3,951,815 and 55-38164 (U.S. Pat. No. 4,005,012), PB Report No. 80-182090 and U.S. Pat. Nos. 4,387,024 and 4,559,139.\nThose prior semipermeable composite membranes mentioned above, however, can not satisfy all of the desired properties required to the membranes, such as high solute rejection against both organic and inorganic materials dissolved in water, high water flux rate, durability such as heat resistance, chemical resistance, pressure resistance, chlorine resistance and hydrogen peroxide resistance.\nThe reverse osmosis is generally expected to be applied to various fields, such as, desalination represented by the conversion of sea water or brackish water into portable water, concentration or recovery of valuable materials, production of ultrapure water and separation of organic liquid mixture. Accordingly, the desired specific properties required to the reverse osmosis membrane would vary multifariously depending on its use and application. The important properties among them are selective permeability (solute rejection: Rej (%), water flux: Flux (m.sup.3 /m.sup.2.day)), durability (m-value relating to the change of flux with the lapse of time), pressure resistance, heat resistance, oxidation resistance, workable pH range, resistance to microorganism, durability for chemicals and solvent, storage stability, etc. The m-value mentioned above is represented by the formula log (Qt.sub.2 /Qt.sub.1)/log (t.sub.2 /t.sub.1) where t.sub.1 is an initial time and t.sub.2 is a time after a predetermined period and Qt.sub.1 and QT.sub.2 represent the flux at the time t.sub.1 and the time t.sub.2, respectively.\nHowever, no asymmetric nor semipermeable composite membranes fulfilling all of these requirements has been developed at present. For example, the prior composite reverse osmosis membrane having an ultrathin membrane of cross-linked polyamide is said to show significant reduction of the flux with the lapse of operation period. While the desirable m-value for practical use of the reverse osmosis membrane is usually around -0.02, which corresponds to about 20% reduction of the flux during the operation of three years, any reverse osmosis membrane which has the m-value far from the above can not be considered useful from the viewpoint of efficiency and economization, even though its initial characteristics and properties are excellent and acceptable.\nWhile on the other hand, the practical level of the solute rejection against organic material is usually required to be higher than 90%, representatively, against neutral, basic and acidic molecules such as isopropyl alcohol, ethylenediamine and acetic acid, no reverse osmosis membrane capable of fulfilling this requirement has been developed.\nIn addition, although some of the prior reverse osmosis membranes made of polyamide or cross-linked polyamide had been considered as chlorine resistant, their water flux is, in fact, usually reduced by about 20 to 30% or more after they are brought into contact with chlorine (refer to Journal of Applied Polymer Science, Vol. 29, 3369 (1984)), and the loss of the rejection and increase in the flux rate occur in an unrecoverable manner possibly due to the decomposition of polyamide upon contact with chlorine for a long period of operation or at high chlorine concentration.\nA practically attractive reverese osmosis membrane should have the salt rejection of 99% or higher and water flux of 0.8 m.sup.3 /m.sup.2.day or greater under the condition of the applied pressure of 56 kg/cm.sup.2 and the feed saline water of 35,000 ppm in the case of sea water desalination, or the salt rejection of 99% or higher and water flux of 0.7 to 1.3 m.sup.3 /m.sup.2.day or greater under the conditions of the applied pressure of 15 to 30 kg/cm.sup.2 and the feed saline water of 1500 ppm in the case of brackish water desalination, in addition to a small reduction of water flux with the lapse of operation period, an improved rejection against organic materials and an oxidation resistance. The reverse osmosis membranes having been developed so far can only partially satisfy these requirements, and accordingly are subjected to various restrictions when put into practical use.\nThe object of the invention is to provide a semipermeable composite membrane of which ultra-thin membrane is made of specific cross-linked aromatic polyamide, having superior hydrophilicity, and high cross-linking density to improve especially the solute rejection against both inorganic and organic materials and water flux as well as the mechanical durability of the ultra-thin membrane.\nOther objects and advantages of the invention will be apparent from the descriptions hereinafter."}
{"text": "1. Field of the Invention\nThe present invention relates generally to soil openers that use a disc assembly to open the soil prior to injection of a fertilizer or seed into the ground. More specifically, the present invention concerns a sectored disc assembly for use in such a soil opener, a removable disc sector of a sectored disc assembly, as well as a method of replacing a disc sector of a sectored disc assembly.\n2. Discussion of the Prior Art\nThose of ordinary skill in the art will appreciate that farmers often use no-till planting techniques to produce all types of agricultural crops. Such no-till techniques minimize the disturbance to the soil and leaves the stubble, or organic matter from the previous crop, standing in the field. This, in turn, reduces water runoff in the field, thereby greatly reducing erosion of the top soil.\nTypically, no-till planting involves using a rotatable disc soil opener to cut a furrow in the soil as it is pulled across a field, creating a small disturbed soil zone. Fertilizer, seeds, or both, are then injected into this disturbed soil zone, after which the furrow is closed. The closing of the furrow can be accomplished by naturally allowing the disturbed soil to flow back into the furrow, or by following the rotatable disc soil opener with a closing wheel designed to push soil into the disturbed soil zone and close the furrow. Often, the opening of the furrow, formation of the seed bed, injection of fertilizer, seed, or both, and closing of the furrow are accomplished in a single pass with a soil opener that includes a rotatable disc, an injector foot, and a closing wheel.\nConventional rotatable discs used in soil openers have been formed, most typically out of carbon steel, as single, annular bodies for rotation about a spindle on the soil opener. While this unitary construction has been satisfactory in some respects, such solid discs are very heavy components that can be difficult to handle and are often expensive to produce. As the rotatable disc component of a soil opener directly contacts the soil as the implement moves through a field, the disc is subject to considerable stress and is exposed to damage. Damage to a rotatable disc, often the result of transport, inexperienced or poor operators, or hitting a railroad track or rocks in the soil, requires replacement of the entire disc. Such disc replacement is expensive and time consuming, as the entire heavy disc must be removed from the soil opener, typically involving tedious disassembly of the entire supporting frame, and only then can an entire new disc be installed.\nThe cost associated with whole disc replacement can be disproportional to the amount of damage to the disc, such as when a single portion of the edge of the disc is damaged from hitting a rock. Additionally, the replacement of the whole disc results in considerable down time of the implement, as heavy components must be elevated out of the soil to facilitate removal of the disc to be replaced and installation of a new disc. The time requirement associated with such replacement adds to the cost and inconvenience of the periodic and necessary change out of rotatable discs."}
{"text": "The present invention relates generally to ventilation devices and methods and more particularly to a ventilation device suitable for mounting with a vehicle window for ventilation and exhaust of interior air and/or associated airborne impurities or contaminants.\nThere are two common methods that people use to ventilate a vehicle's interior. One method is to open the vehicle's windows when the vehicle is moving or parked. This method allows smoke, pet smells, and other foul odors to be ventilated to the outside atmosphere, while simultaneously allowing the fresh, outside air to enter the vehicle's cabin. There are, however, some disadvantages to opening the vehicle's windows for ventilation. First, the outside air may not be fresh and clean at all times. For example, the outside air may include rain, snow, pollen, bugs or any number of other outside contaminants, which will also enter the vehicle's cabin. Second, while the vehicle is in a parked position and left unattended, having the windows down may result in the vehicle's theft or vandalism. Third, if the windows are down while the vehicle is in motion, the vehicle's aerodynamic qualities will be affected significantly. Approximately 10% of the vehicle's fuel economy will be lost due to drag forces. Finally, if the windows are down while the vehicle is in motion, a significant noise is produced which will cause discomfort to the ordinary passenger. The passengers will be required to converse at a higher decibel and also be required to turn vehicle audio equipment to a louder setting.\nAnother method commonly used to ventilate the vehicle is to use the ventilation feature on the vehicle's control panel. This feature allows fresh air at a particular temperature to enter the vehicle's cabin. While the fan functions adequately in most situations, there are some scenarios where the fan still leaves behind undesirable odors. For example, some odors which have penetrated the seat's fabric are not capable of being removed by the fans. Instead, the high velocity air entering the vehicle's window, while the vehicle is moving at high speeds, is more suitable for removing these types of odors.\nVarious solutions have been proposed to resolve the problem of safe vehicle ventilation. One such solution proposes using a panel mountable between a window frame and a window. Although prior art devices allow air to exchange between the vehicle's exterior and its interior, only one of these prior art devices forcefully removes air from the vehicle's interior by creating a pressure differential between the vehicle's interior and its exterior. Forceful air removal increases the speed with which the vehicle's interior is ventilated and the obnoxious odors are removed. Several drawbacks exist in the current art which include, but are not limited to, requiring the mounting of a device between the window and the window frame, the window not being capable of fully closing, inefficient pressure differential along the device, and water leakage possibilities around the device.\nWhat has been lacking, however, until the present invention, and what the industry long has sought, is a cabin exhaust device that has flexibility in its mounting, allows the windows to fully close when exhaust is not needed, creates a more efficient pressure differential throughout the device's length and prevents water leakage around the device.\nIt is, therefore, an object of the present invention to disclose and claim a cabin exhaust device for use in vehicles with retractable windows.\nIt is another object of the present invention to provide a cabin exhaust device that facilitates the creation of a vacuum, or pressure differential, between the vehicle's interior and its exterior and along the entire length of the cabin exhaust device.\nIt is another object of the present invention to provide a cabin exhaust device that allows the windows to fully close while the cabin exhaust device is mounted.\nIt is another object of the present invention to provide a cabin exhaust device that vents cabin heat and removes odors when the vehicle is either parked or in motion.\nIt is yet another object of the present invention to provide a cabin exhaust device which keeps out dust, pollen, allergens, and other similar pollutants.\nIt is a further object of the instant invention to provide a cabin exhaust device which is easy to position and mount.\nIt is yet another object of the present invention to provide a cabin exhaust device which can be mounted on the outside of the door frame or in an opening between the window panel and the window frame.\nIt is still a further object of the present invention to provide a cabin exhaust device which is inexpensive to manufacture.\nIt is yet another object of the present invention to provide a cabin exhaust device which helps lower the temperature inside the cabin when the vehicle is stationary, thereby preventing the deterioration of adhesives within the vehicle, extend the life of plastics and other materials within the vehicle, and provide environmental safety to the vehicle's passengers by eliminating plastic and glue fumes as well as other contaminants that accumulate in the vehicle's cabin.\nIt is yet another object of the present invention to provide a cabin exhaust device which provides cabin ventilation in a fuel efficient manner.\nIt is yet another object of the present invention to provide a cabin exhaust device that prevents rain from entering the vehicle's cabin.\nIt will become apparent to one skilled in the art that the claimed subject matter as a whole, including the structure of the apparatus, and the cooperation of the elements of the apparatus, combine to result in the unexpected advantages and utilities of the present invention. The advantages and objects of the present invention and features of such a cabin exhaust device will become apparent to those skilled in the art when read in conjunction with the accompanying description, drawing figures, and appended claims."}
{"text": "The present invention relates to universal data output and, in particular, to providing a new data output method and a new raster image process for information apparatuses and output devices.\nAs described herein, information apparatuses refer generally to computing devices, which include both stationary computers and mobile computing devices (pervasive devices). Examples of such information apparatuses include, without limitation, desktop computers, laptop computers, networked computers, palmtop computers (hand-held computers), personal digital assistants (PDAs), Internet enabled mobile phones, smart phones, pagers, digital capturing devices (e.g., digital cameras and video cameras), Internet appliances, e-books, information pads, and digital or web pads. Output devices may include, without limitation, fax machines, printers, copiers, image and/or video display devices (e.g., televisions, monitors and projectors), and audio output devices.\nFor simplicity and convenience, hereafter, the following descriptions may refer to an output device as a printer and an output process as printing. However, it should be understood that the term printer and printing used in the discussion of present invention refer to one embodiment used as a specific example to simplify the description of the invention. The references to printer and printing used here are intended to be applied or extended to the larger scope and definition of output devices and should not be construed as restricting the scope and practice of present invention.\nFueled by an ever-increasing bandwidth, processing power, wireless mobile devices, and wireless software applications, millions of users are or will be creating, downloading, and transmitting content and information using their pervasive or mobile computing devices. As a result, there is a need to allow users to conveniently output content and information from their pervasive computing devices to any output device. As an example, people need to directly and conveniently output from their pervasive information apparatus, without depending on synchronizing with a stationary computer (e.g., desktop personal computer) for printing.\nTo illustrate, a mobile worker at an airport receiving e-mail in his hand-held computer may want to walk up to a nearby printer or fax machine to have his e-mail printed. In addition, the mobile worker may also want to print a copy of his to-do list, appointment book, business card, and his flight schedule from his mobile device. As another example, a user visiting an e-commerce site using his mobile device may want to print out transaction confirmation. In still another example, a user who takes a picture with a digital camera may want to easily print it out to a nearby printer. In any of the above cases, the mobile user may want to simply walk up to a printer and conveniently print a file (word processing document, PDF, HTML etc) that is stored on the mobile device or downloaded from a network (e.g., Internet, corporate network).\nConventionally, an output device (e.g., a printer) is connected to an information apparatus via a wired connection such as a cable line. A wireless connection is also possible by using, for example, radio communication or infrared communication. Regardless of wired or wireless connection, a user must first install in the information apparatus an output device driver (e.g., printer driver in the case the output device is a printer) corresponding to a particular output device model and make. Using a device-dependent or specific driver, the information apparatus may process output content or digital document into a specific output device's input requirements (e.g., printer input requirements). The output device's input requirements correspond to the type of input that the output device (e.g., a printer) understands. For example, a printer's input requirement may include printer specific input format (e.g., one or more of an image, graphics or text format or language). Therefore, an output data (or print data in the case the output device is a printer) herein refers to data that is acceptable for input to an associated output device. Examples of input requirements may include, without limitation, audio format, video format, file format, data format, encoding, language (e.g., page description language, markup language etc), instructions, protocols or data that can be understood or used by a particular output device make and model.\nInput requirements may be based on proprietary or published standards or a combination of the two. An output device's input requirements are, therefore, in general, device dependent. Different output device models may have their own input requirements specified, designed or adopted by the output device manufacturer (e.g., the printer manufacturer) according to a specification for optimal operation. Consequently, different output devices usually require use of specific output device drivers (e.g., printer drivers) for accurate output (e.g., printing). Sometimes, instead of using a device driver (e.g., printer driver), the device driving feature may be included as part of an application software.\nInstallation of a device driver (e.g., printer driver) or application may be accomplished by, for example, manual installation using a CD or floppy disk supplied by the printer manufacturer. Or alternatively, a user may be able to download a particular driver or application from a network. For a home or office user, this installation process may take anywhere from several minutes to several hours depending on the type of driver and user's sophistication level with computing devices and networks. Even with plug-and-play driver installation, the user is still required to execute a multi-step process for each printer or output device.\nThis installation and configuration process adds a degree of complexity and work to end-users who may otherwise spend their time doing other productive or enjoyable work. Moreover, many unsophisticated users may be discouraged from adding new peripherals (e.g., printers, scanners, etc.) to their home computers or networks to avoid the inconvenience of installation and configuration. It is therefore desirable that an information apparatus can output to more than one output device without the inconvenience of installing multiple dedicated device dependent drivers.\nIn addition, conventional output or printing methods may pose significantly higher challenges and difficulties for mobile device users than for home and office users. The requirement for pre-installation of a device-dependent driver diminishes the benefit and concept of mobile (pervasive) computing and output. For example, a mobile user may want to print or output e-mail, PowerPoint® presentation documents, web pages, or other documents at an airport, gas station, convenience store, kiosk, hotel, conference room, office, home, etc. It is highly unlikely that the user would find at any of these locations a printer of the same make and model as is at the user's base station. As a consequence, under the conventional printing method, the user would have to install and configure a printer driver each time at each such remote location before printing. It is usually not a viable option given the hundreds, or even thousands of printer models in use, and the limited storage, memory space, and processing power of the information apparatus.\nMoreover, the user may not want to be bothered with looking for a driver or downloading it and installing it just to print out or display one page of email at the airport. This is certainly an undesirable and discouraging process to promote pervasive or mobile computing. Therefore, a more convenient printing method is needed in support of the pervasive computing paradigm where a user can simply walk up to an output device (e.g., printer or display device) and easily output a digital document without having to install or pre-install a particular output device driver (e.g., printer driver).\nAnother challenge for mobile users is that many mobile information apparatuses have limited memory space, processing capacity and power. These limitations are more apparent for small and low-cost mobile devices including, for example, PDAs, mobile phones, screen phones, pagers, e-books, Internet Pads, Internet appliances etc. Limited memory space poses difficulties in installing and running large or complex printer or device drivers, not to mention multiple drivers for a variety of printers and output devices. Slow processing speed and limited power supply create difficulties driving an output device. For example, processing or converting a digital document into output data by a small mobile information apparatus may be so slow that it is not suitable for productive output. Intensive processing may also drain or consume power or battery resources. Therefore, a method is needed so that a small mobile device, with limited processing capabilities, can still reasonably output content to various output devices.\nTo output or render content (e.g. digital document) to an output device, a raster image processing (RIP) operation on the content is usually required. RIP operation can be computationally intensive and may include (1) a rasterization operation, (2) a color space conversion, and (3) a halftoning operation. RIP may also include other operations such as scaling, segmentation, color matching, color correction, GCR (Grey component replacement), Black generation, image enhancement compression/decompression, encoding/decoding, encryption/decryption GCR, image enhancement among others.\nRasterization operation in RIP involves converting objects and descriptions (e.g. graphics, text etc) included in the content into an image form suitable for output. Rasterization may include additional operations such as scaling and interpolation operations for matching a specific output size and resolution. Color space conversion in RIP includes converting an input color space description into a suitable color space required for rendering at an output device (e.g. RGB to CMYK conversion). Digital halftoning is an imaging technique for rendering continuous tone images using fewer luminance and chrominance levels. Halftoning operations such as error diffusion can be computationally intensive and are included when the output device's bit depth (e.g. bits per pixel) is smaller than the input raster image bit depth.\nConventionally, RIP operations are included either in an information apparatus, or as part of an output device or output system (e.g. in a printer controller). FIG. 1A illustrates a flow diagram of a conventional data output method 102 in which RIP 110 is implemented in the information apparatus. Output devices that do not include a printer controller to perform complex RIP operations, such as a lower-cost, lower speed inkjet printer, normally employ data output method 102. In data output method 102, an information apparatus obtains content (e.g. a digital document) in step 100 for rendering or output at an output device. The information apparatus may includes an application (e.g. device driver), which implements RIP operation 110. The information apparatus generates an output data in step 120 and transmits the output data to the output device in step 130 for rendering. The output data relating to the content is in an acceptable form (e.g. in an appropriate output size and resolution) to the output engine (e.g. display engine, printer engine etc.) included in the output device. The output data in a conventional output method 102 is usually device dependent.\nOne drawback for the data output method 102 of FIG. 1A is that the information apparatus performs most if not the entire raster image processing operations 110 required for output. The RIP operations may require intensive computation. Many information apparatus such as mobile information device might have insufficient computing power and/or memory to carry out at an acceptable speed the RIP operations 110 required in an output process.\nAnother drawback for the conventional data output method 102 of FIG. 1A is that the generated output data is device dependent and therefore is typically not very portable to other output devices. As a result, the information apparatus may need to install multiple applications or device drivers for multiple output devices, which may further complicate its feasibility for use in information apparatuses with limited memory, storage and processing power.\nFIG. 1B illustrates a flow diagram of another conventional data output method 104 in which the RIP is implemented in an output device. An example of an output device that implements process 104 is a high-speed laser printer which includes a printer controller for performing RIP operations and an output engine (e.g. printer engine) for rendering content. Printer controller may be internally installed or externally connected to an output device (printer in this example). In data output method 104, an information apparatus obtains content for output in step 100 and generates in step 160 an output data or print data for transmitting to the output device in step 170. Print data includes information related to the content and is usually encoded in a page description language (PDL) such as PostScript and PCL etc. In step 180, the printer receives the output data or print data (in a PDL). In step 190, a printer controller included in the printer interprets the PDL, performs RIP operations, and generates a printer-engine print data that is in a form acceptable to the printer engine (e.g. a raster image in an appropriate output size, bit depth, color space and resolution). In step 150 the printer engine renders the content with the printer-engine print data.\nIt will be understood that a reference to print data or output data including a language, such as PDL, should be interpreted as meaning that the print data or output data is encoded using that language. Correspondingly, a reference to a data output process generating a language, such as PDL, should be interpreted as meaning that the data output process encodes data using that language.\nThere are many drawbacks in the conventional data output method 104 shown in FIG. 1B. These drawbacks are especially apparent for mobile computing devices with limited processing power and memory. One such drawback is that the output data or print data, which include a page description language (PDL) such as PostScript or PCL, can be very complex. Generating complex PDL may increase memory and processing requirements for an information apparatus. Furthermore, interpreting, decoding and then raster image processing complex PDL can increase computation, decrease printing speed, and increase the cost of the output device or its printer controller.\nAnother drawback is that the output data that includes PDL can creates a very large file size that would increase memory and storage requirements for the information apparatus, the output device and/or the printer controller etc. Large file size may also increase the bandwidth required in the communication link between the information apparatus and the output device.\nFinally, to rasterize text in an output device, a printer controller may need to include multiple fonts. When a special font or international characters is not included or missing in the printer controller, the rendering or output can potentially become inaccurate or inconsistent."}
{"text": "A cis-4-fluoro-L-proline derivative can be prepared from a trans-4-hydroxy-L-proline derivative through fluorination of its hydroxyl group at the 4-position. However, conventional fluorination techniques have posed various problems, as explained below.\nThe first method for fluorination is a method which uses diethylaminosulfur trifluoride (DAST). This method provides a target product in good yield, but is unsuitable for industrial use because DAST is highly toxic and less heat stable, as well as being explosive and expensive (Luc Demange et al., Tetrahedron Lett., 39, 1169 (1998)).\nThe second method is a method in which the hydroxyl group at the 4-position is converted into a leaving group and then into a fluoro group (G. Giardina et al., Synlett, 1, 55 (1995)). This method is of limited practical use because the yield of a target product is reduced as a result of olefin generation caused by β-elimination of the leaving group.\nThe third method is a method which uses N,N-diethyl-N-(1,1,2,3,3,3-hexafluoropropyl)amine (hereinafter referred to as “Ishikawa reagent”). This method is considered to be excellent in that it has a reduced risk of explosion, is available at low cost and enables the introduction of a fluoro group in a single step. However, there is a problem of reduced yield of a target compound because hydrogen fluoride generated during the reaction causes not only erosion of reaction vessels, but also decomposition of starting materials and reaction products. Moreover, particularly in the case of using a urethane-type protecting group which is widely used as a protecting group for an amino group in an α-amino acid, this protecting group will be easily decomposed by the action of hydrogen fluoride and hence requires the limited reaction conditions (e.g., at low temperature over a long period of time), thus failing to produce satisfactory results in terms of yield, etc."}
{"text": "This invention relates generally to golfing, and more particularly to golf ball retrieving and extracting devices that fit or attach to golf clubs.\nNormally, at the completion of each hole golfers must bend over to remove a golf ball from the cup. Also, oftentimes during play golfers may have to bend over numerous times to remove golf balls from hazards, such as grass or water. Although bending over is usually not a problem for the young golfer, it can be a strenuous and virtually impossible task for the elderly, anyone with back problems, arthritis, or a handicapped golfer. Thus, there is a need for a device that will extract golf balls from such locations without requiring the golfer to bend over.\nFurthermore, such a ball extracting device must be convenient and easy to use as usually a golfer carries only the putter in hand at the end of the hole or, if not a that hole, the golfer carries another club and is separated from the golf bag. Therefore, it is necessary to have such a device that is built into the putter or club which, at the same time, does not effect the playing characteristics of the putter or club.\nThe prior art contains many devices designed to fulfill the above needs, but none does it like the present invention. The most pertinent patents concerning golf ball retrievers or extractors that attach to the handle-grip end of a golf club are as follows:\n______________________________________ U.S. Pat. No.: Inventor: Date of Patent: ______________________________________ 5,004,240 Tsukamoto April 2, 1991 4,787,632 Nigrelli, et al. Nov. 29, 1988 3,318,628 White May 9, 1967 3,401,970 Russell Sept. 17, 1968 3,462,184 Russell Aug. 19, 1969 2,833,584 McEvoy May 6, 1958 2,801,875 McEvoy Aug. 6, 1957 802,264 Brown Oct. 17, 1905 3,698,720 Gudmundsen Oct. 17, 1972 3,273,927 Carlson Sept. 20, 1966 3,239,264 DuPont March 25, 1964 4,687,204 Lempio Aug. 18, 1987 1,380,526 Carpenter June 7, 1921 1,674,294 O'Rourke June 19, 1928 2,523,942 Ciambriello Sept. 26, 1950 2,750,222 Borah June 12, 1956 ______________________________________\nThe Tsukamoto patent shows an extractor will ball-retrieving prongs built into the grip handle. Unfortunately, the latter device requires that the golf club handle be modified so that a spring and other mechanical components can be placed therein to operate the retrieving prongs. The Nigrelli, et al. patent shows claw-clamp type retrievers that attach to the end of a club handle. The White patent shows another prong-like retriever built into the handle grip operable by a push button extension. The two Russel patents show other prong-like retrievers that attach to the end of a handle or which are built into the handle grip extension. The two McEvoy patents also show palm retrievers that are removably attachable to the handle by different means. The remainder of the above patents show various golf ball retrievers that fit by a collar over the end of the golf club or are incorporated within the handle grip itself.\nUnlike the golf ball retrievers or extractors in the prior art that are attached to the handle end of the golf club, the present invention is an integral part of the grip itself and thus requires no modification of an existing golf club. Furthermore, unlike inventions in the prior art the present invention has no moving parts. The present invention discloses an extractor which is a circular module having an opening on the top end which holds a ball when it is pushed downward on the ball. The present extractor has one or more finger holes around the extractor to enable the ball to be released from the extractor by use of the golfer's fingers. One embodiment of the present invention has an entire grip that fits over the golf club shaft and replaces the preexisting grip on the golf club. Another embodiment is designed to fit over a preexisting grip. Finally, a third embodiment attaches just over the tip of the grip itself. All of the embodiments require no modification to existing golf clubs and have no impact upon putting or playing performance."}
{"text": "Description of the Prior Art\nThis invention, in its preferred form, relates to welding fixture for welding sleeves to a fuel grid and to a method for assembling and locating the sleeves relative to the grid. More particularly, the invention relates to the apparatus, and the attendant method, for locating a grid on a fixture, and for locating sleeves on the grid, through use of the fixture. As herein used, the word \"fixture\" includes both a plate having abutments and locating pins, and independent pegs.\nThere have been provided nuclear fuel bundle assemblies which include a matrix of nuclear fuel rods which are arrayed in rows and columns and held in such configuration by a plurality of fuel rod grids. These grids are formed of straps which are slotted so that the straps mate, one strap entering into the slot of the outer strap at each point of intersection. Control rods are included at selected positions within the array of nuclear fuel rods. These control rods are received in guide sleeves which are welded to the grid. More particularly, there has been provided in the past a construction in which the straps forming the grid had edges which were substantially in a common plane, and a plurality of guide sleeves were provided, each seated on four straps defining a common cell. The outside diameter of the cylindrical guide sleeve was such as to permit the base of the cylindrical guide sleeve to rest upon the edges of four straps forming a single cell, it being understood that guide sleeves were provided only at spaced cells, and not at adjacent cells.\nTo assemble the guide sleeve to the grid, the guide sleeves were butt welded onto the straps forming the particular cell at the upper edge of each of the four straps. This procedure required that the grid, after assembly, be machined to insure that the edges of the four straps forming the particular cell to which the cylindrical sleeve was to be welded lay as closely as possible in a common plane, and this required a fabrication step known as \"spot facing\". The \"spot facing\" step was time consuming and costly, requiring the use of a suitable machine and fixtures, as well as the expenditure of a workman's time. The requirement for relatively great precision in the location of the strap edges is based on the fact that if a gap greater than 0.005 inches remains between the sleeve and the edge of the grid strap, the weld which is effected may well be defective in that the fused metal will not extend entirely beneath that portion of the base of the cylindrical sleeve which is intended to engage the upper edge of the grip strap; the occurrence of a gap will result in either lack of welding or no welding, either of which is unsatisfactory.\nWelding fixtures and methods have been heretofore provided, involving the provision of so called \"chills\" generally in the nature of pins which extended into cells and which positioned sleeves on the grid, for welding. These fixtures have not proven satisfactory, particularly in connection with the utilization of laser welding for welding the sleeves to the grid straps."}
{"text": "Physical media can be provided with a purchase of a computer system. The physical media can allow a user to restore the computer system to the factory defaults if there is a problem with the computer system. The physical media can also be used to perform data recovery. However, physical media can be expensive and unreliable. The physical media can be faulty due to a manufacturing process problem or may not be supported by the disk drive used to read the physical media. Furthermore, the physical media is a fixed size which limits the amount of information that can be stored in the physical media. Moreover, the physical media is read-only so it cannot be updated or upgraded by the user. For example, if a user upgrades the operating system on the computer system and needs to restore the computer system, the physical media would restore the computer system to the originally installed operating system (factory default), rather than the upgraded version. If a user wanted to restore to the updated or upgraded version of the operating system, the user would be required to obtain a new physical media corresponding to the upgraded or updated version of the operating system or download and install the updates."}
{"text": "In recent years, with gradual exhaustion of non-renewable resources, such as oil, and increasingly severe environmental damages, more and more attentions are paid on new energy vehicles. As the key role of the new energy vehicle, power battery has obtained increasingly higher attention nowadays.\nHowever, during running of the power battery, there are many uncontrollable and extremely bad situations which may cause short circuit inside the battery. The short circuit may produce a lot of heat, which may damage the battery or harm the users, thus reduce the safety and reliability of the power battery. Accordingly, a safe and reliable battery is in need in this field."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, and more particularly, to a semiconductor device that can have power consumption reduced.\n2. Description of the Background Art\nSemiconductor devices incorporated in various equipments have the scale of integration increased in order to reduce the size of the device and to integrate much more powerful logic. Increase in the integration density will result in a great number of elements operating inside to raise the heat. Therefore, reducing power consumption is an important factor. For example, in a DRAM (Dynamic Random Access Memory), there is a great demand for reducing the power consumption as a result of increase in the number of elements according to increase in the storage capacity.\nA DRAM as a conventional semiconductor memory device will be described in detail hereinafter. A DRAM includes an intermediate potential generation circuit for generating precharge potential for a bit line, a timer circuit for carrying out a self refresh operation, and an internal high voltage circuit for generating high potential to be provided to a word line drive circuit.\nAn intermediate potential generation circuit will first be described. An example of an intermediate potential generation circuit is disclosed in IEEE Journal of Solid-State Circuit, Vol. SC-22, No. 5, October 1987, pp. 861-867. FIG. 25 is a circuit diagram showing a structure of a conventional intermediate potential generation circuit thereof.\nReferring to FIG. 25, an intermediate potential generation circuit includes transistors Q101-Q103 which are n type MOSFETS, transistors Q104-Q106 which are p type MOSFETs, and resistors R101-R104.\nFIG. 26 schematically shows a structure of the intermediate potential generation circuit of FIG. 25 on a p type substrate. Referring to FIG. 26, the intermediate potential generation circuit includes a p type substrate 111, an n type well 112, transistors Q11-Q106, and resistors R101-R104. In FIG. 26, components corresponding to those of FIG. 25 have the same reference characters denoted.\nAn operation of the intermediate potential generation circuit will be described hereinafter with reference to FIGS. 25 and 26.\nThe resistance of resistors R101 and R102 equal each other. Also, the resistance of resistors R103 and R104 equal each other. The resistance of resistors R101-R104 is several m.OMEGA., which is high resistance. Therefore, the current flowing in transistors Q101, Q102, Q104 and Q105 is reduced, and these transistors conduct lightly. Therefore, the gate-source potential of transistor Q101, Q102, Q104 and Q105 is equal to the threshold voltage of each transistor.\nAccording to the above-described structure, the potential of nodes N1 and N3 is approximately V.sub.CC /2 (V.sub.CC is the power supply voltage). Therefore, the potential of node N2 becomes V.sub.CC /2+V.sub.TH101 (V.sub.TH101 is the threshold voltage of transistor Q101), and the potential of node N4 is approximately V.sub.CC /2-.vertline.V.sub.TH105 .vertline. (V.sub.TH105 is the threshold voltage of transistor Q105). When the potential of an output signal V.sub.sg is lower than V.sub.CC /2+V.sub.TH101 -V.sub.TH103 (V.sub.TH103 is the threshold voltage of transistor Q103), transistor Q103 conducts, whereby the potential of output signal V.sub.sg rises. When the potential of output signal V.sub.sg is higher than V.sub.CC /2-.vertline.V.sub.TH105 .vertline.+.vertline.V.sub.TH106 .vertline. (V.sub.TH106 is the threshold voltage of transistor Q106), transistor Q106 conducts, whereby the potential of output signal V.sub.sg falls. By the above-described operation, the potential of output signal V.sub.sg becomes approximately V.sub.CC /2.\nA timer circuit for a self refresh operation will be described hereinafter. A refresh operation must be carried out periodically since a DRAM is a volatile memory. Lengthening the period of a refresh operation will reduce power consumption thereof, to allow reduction of power consumption in the device. In a conventional timer circuit, a refresh operation is carried out when the potential held in a memory cell becomes lower than a predetermined level. An example of such a timer circuit is disclosed in IEEE Journal of Solid-State Circuits, Vol. 26, No. 11, November 1991, pp. 1556-1562. FIG. 27 shows a structure of this conventional timer circuit.\nReferring to FIG. 27, a timer circuit includes a differential amplifier 121, an S-R flipflop 122, a delay circuit 123, a transistor Q111 which is an n type MOSFET, a capacitor 124 of a memory cell, and an n type diffusion layer 125.\nFIGS. 28A(1) to 28B(B) are timing chart showing the operation of the timer circuit of FIG. 27.\nAn operation of the timer circuit will be described hereinafter with reference to FIGS. 27 and 28A(1) to 28B(B). When the potential V.sub.N in capacitor 124 becomes lower than a reference potential V.sub.REF at time t.sub.1, S-R flipflop 122 is set to render the level of an output signal .phi..sub.E to a H level (logical high). Output signal .phi..sub.E of S-R flipflop 122 is delayed for a predetermined time, and then applied to a reset terminal R of S-R flipflop 122. As a result, a reset signal R attains a H level. This causes S-R flipflop 122 to be reset, whereby output signal .phi..sub.E attains a L level (logical low). A refresh operation is carried while output signal .phi..sub.E attains a H level, whereby transistor Q111 attains a conductive state, and the potential of capacitor 124 of a memory cell is maintained at V.sub.CC. Then, when output signal .phi..sub.E attains a L level, transistor Q111 is rendered non-conductive, whereby the holding voltage V.sub.N of capacitor 124 is gradually reduced by leakage current. When holding voltage V.sub.N of capacitor 124 becomes lower than reference voltage V.sub.REF, an operation similar to that of the above-described operation is repeated. Thus, a refresh operation is carried out at a predetermined period.\nAn internal high voltage circuit will be described. FIG. 29 is a block diagram showing a structure of a conventional internal high voltage circuit. Referring to FIG. 29, an internal high voltage circuit includes a first detector 132, a second detector 132, a third detector 133, a first oscillator 134, a second oscillator 135, a small pump 136, a large pump 137, a RAS pump 138, and an AND gate G101 and an inverter G102.\nWhen high voltage V.sub.PP supplied to a word line driver 139 becomes lower than a predetermined potential, first detector 131 provides an output signal .phi..sub.E1 of a H level to first oscillator 134. First oscillator 134 oscillates while output signal .phi..sub.E1, attains a H level, and provides an oscillation signal to small pump 136. Small pump 136 responds to this oscillation signal to provide high voltage V.sub.PP to word line driver 139 at a standby state.\nWhen the high voltage supplied to word line driver 139 becomes lower than a predetermined potential, second detector 132 provides an output signal .phi..sub.E2 of H level to second oscillator 135. Second oscillator 135 oscillates when output signal .phi..sub.E2 attains a H level, and provides an oscillation signal to large pump 137. Large pump 137 responds to this oscillation signal to rapidly increase high voltage V.sub.PP supplied to word line driver 139.\nWhen high voltage V.sub.PP supplied to word line driver 139 becomes lower than a predetermined potential, third detector 133 provides an output signal .phi..sub.E3 of a H level to AND gate G101. AND gate G101 takes the logical product of output signal .phi..sub.E3 and an inverted signal of a row address strobe signal /RAS (\"/\" implies a low-active signal) to provide an output signal to RAS pump 138. AND gate G101 provides an output signal when row address strobe signal /RAS attains a L level, whereby the semiconductor device operates to raise the word line to high voltage V.sub.PP.\nThe first detector shown in FIG. 29 will be described with reference to FIG. 30 showing a circuit diagram thereof.\nReferring to FIG. 30, a first detector includes transistors Q121-Q124 which are p type MOSFETs, and transistors Q125 and Q126 which are n type MOSFETs.\nHigh voltage V.sub.PP provided to the first detector is reduced by a threshold voltage .sub.TH of each transistor, i.e., by 3V.sup.TH, by transistors Q121-Q123. Therefore, an output signal .phi..sub.E1 of a H level is output when high voltage V.sub.PP becomes lower than V.sub.CC +3V.sub.TH.\nThe second detector of FIG. 29 will be described hereinafter with reference to FIG. 31 showing a circuit diagram thereof.\nReferring to FIG. 31, a second detector includes transistors Q131-Q133 which are p type MOSFETs, and transistors Q134 and Q135 which are n type MOSFETs.\nHigh voltage V.sub.PP provided to the second detector is reduced by a threshold voltage V.sub.TH of each transistor, i.e. 2V.sup.TH, by transistors Q131 and Q132. Therefore, the second detector provides an output signal .phi..sub.E2 of a H level when high voltage V.sub.PP becomes lower than V.sub.CC +2V.sub.TH. The third detector of FIG. 29 has a structure similar to that of the second detector of FIG. 31, and also operates in a similar manner thereof.\nThe first oscillator of FIG. 29 will be described hereinafter with reference to FIG. 32 showing a circuit diagram thereof.\nReferring to FIG. 32, a first oscillator includes transistors Q141-Q148 which are p type MOSFETs, and transistors Q149-Q156 which are n type MOSFETs. C101-C103 shown in FIG. 20 are the parasitic capacitance of each portion.\nBecause transistor Q141 has a long channel length, current flowing in transistor Q149 is limited to a current value of I.sub.1. Transistor Q149 and transistors Q150, Q152, Q154 and Q156 form a current mirror, so that current flowing through transistors Q143, Q145, Q147, Q152, Q154 and Q156 is limited to the value of I.sub.1. Therefore, the delay time of each inverter formed by each of these transistors becomes 3C/I.sub.1 where each capacitance of parasitic capacitances C101-C103 is C.\nWhen V.sub.CC /2+V.sub.TH101 -V.sub.TH103 &gt;V.sub.sg =V.sub.CC /2, and V.sub.CC /2-.vertline.V.sub.TH105 .vertline.+.vertline.V.sub.TH106 .vertline.&lt;V.sub.sg =V.sub.CC /2, i.e. V.sub.TH101 &gt;V.sub.H103, and .vertline.V.sub.TH105 .vertline.&gt;.vertline.V.sub.TH106 .vertline. in the intermediate potential generation circuit of FIG. 25, a through current flows in transistors Q103 and Q106 at the time of standby since transistors Q103 and Q106 both conduct when the potential of output signal V.sub.sg is stable at V.sub.CC /2. There was a problem that the power consumption of the device was increased due to this through current.\nIn the timer circuit of FIG. 27, the period of a refresh operation is T.sub.1, at low temperature as shown in FIG. 28(a), and is T.sub.2 at high temperature as shown in FIG. 28(b) because leakage current of capacitor 124 increases at high temperature.\nThe timer circuit shown in FIG. 27 had problems set forth in the following. A phenomenon called soft error is seen in a DRAM. More specifically, .alpha. particles emitted from the package or the like cause the generated electrons to be captured in an n type diffusion layer 125 of a memory cell, whereby information in the memory cell is destroyed. Therefore, soft error easily occurs when the holding voltage V.sub.N becomes not higher than the lowest holding voltage V.sub.REF required for proper operation of a readout circuit of a memory cell by a predetermined value of .DELTA.V. As a result, when the level of holding voltage V.sub.REF is equal at both the high and low temperature, the time period having a high probability of generating soft error is d.sub.1, and d.sub.2 at a low temperature and a high temperature, respectively, as shown in FIGS. 28A and 28B. Therefore, there was a problem that the possibility of soft error occurrence is increased at low temperature.\nIn the first and second detectors shown in FIGS. 30 and 31, a through current is conducted to increase power consumption since all transistors Q124, Q126, Q133 and Q135 become conductive when the level of output signals .phi..sup.E1 and .phi..sub.E2 change.\nIn the third detector shown in FIG. 31, the time required for pulling the potential of the node between transistorS Q132 and Q134 from a H level to a L level is several .mu.s. An operation of a DRAM occurs at the minimum of every 90 ns, for example. Therefore, a word line is driven several ten times during the transition of the third director from an off state to an on state, resulting in reduction of the level of high voltage V.sub.PP of word line driver 139. FIGS. 33(a) to 33(b) are diagrams for describing the change in the level of high voltage V.sub.PP, with respect to output signal .phi..sub.E3 of the third director. It is appreciated from FIG. 33(a) to 33(b) that the level of high voltage V.sub.PP is gradually reduced according to each transition of row address strobe signal /RAS when output signal .phi..sub.E3 attains a L level. Therefore, a conventional third detector is set so that sufficient current is conducted to transistor Q134 in order to rapidly pull down the potential of the node between transistors Q132 and Q134 rapidly to a L level from a H level. Thus, there was a problem that power consumption is increased during standby.\nIn the first oscillator of FIG. 32, the delay time of 3C/I.sub.1 is reduced due to increase of current I.sub.1 flowing in transistor Q141 in response to increase of power supply potential V.sub.CC. This causes the oscillation frequency of the first oscillator to be increased to shorten the operation cycle. Thus, there was a problem that power consumption of the device is increased."}
{"text": "Isoprenoids are compounds derived from the five-carbon molecule, isopentenyl pyrophosphate. Investigators have identified over 29,000 individual isoprenoid compounds, with new ones continuously being discovered. Isoprenoids are often isolated from natural products, such as plants and microorganisms, which use isopentenyl pyrophosphate as a basic building block to form relatively complex structures. Vital to living organisms, isoprenoids serve to maintain cellular fluidity and electron transport, as well as function as natural pesticides, to name just a few of their roles in vivo. Furthermore, the pharmaceutical and chemical communities use isoprenoids as pharmaceuticals, nutriceuticals, flavoring agents, and agricultural pest control agents. Given their importance in biological systems and usefulness in a broad range of applications, isoprenoids have been the focus of much attention by scientists.\nConventional means for obtaining isoprenoids include extraction from biological materials (e.g., plants, microbes, and animals) and partial or total organic synthesis in the laboratory. Such means, however, have generally proven to be unsatisfactory. For example, organic synthesis is usually complex since several steps are required to obtain the desired product. Furthermore, these steps often involve the use of toxic solvents, which require special handling and disposal. Extraction of isoprenoids from biological materials may also require toxic solvents. In addition, extraction and purification methods usually provide a low yield of the desired isoprenoid, as biological materials typically contain only small quantities of these compounds. Unfortunately, the difficulty involved in obtaining relatively large amounts of isoprenoids has limited their practical use. In fact, the lack of readily available methods by which to obtain certain isoprenoids has slowed down the progression of drug candidates through clinical trials. Furthermore, once an isoprenoid drug candidate has passed the usual regulatory scrutiny, the actual synthesis of the isoprenoid drug may not lend itself to a commercial scale.\nAs a solution to such problems, researchers have looked to biosynthetic production of isoprenoids, Some success has been obtained in the identification and cloning of the genes involved in isoprenoid biosynthesis. For example, U.S. Pat. No. 6,291,745 to Meyer et al. describes the production of limonene and other metabolites in plants. Although many of the genes involved in isoprenoid biosynthesis may be expressed in functional form in Escherichia coli and other microorganisms, yields remain relatively low as a result of minimal amounts of precursors, namely isopentenyl pyrophosphate.\nIn an effort to address the lack of isopentenyl pyrophosphate, some investigators have attempted to increase isopentenyl pyrophosphate production. Croteau et al. describe in U.S. Pat. No. 6,190,895 the nucleic acid sequences that code for the expression of 1-deoxyxylulose-5-phosphate synthase, an enzyme used in one biological pathway for the synthesis of isopentenyl pyrophosphate. Low yields of isopentenyl pyrophosphate remain, however, since several more enzymes are needed to catalyze other steps in this isopentenyl pyrophosphate biosynthetic pathway. Further, the reference does not address an alternative pathway for isopentenyl pyrophosphate biosynthesis, namely the mevalonate pathway.\nThus, the current invention is directed toward solving these and other disadvantages in the art by increasing the typically low yields associated with conventional synthesis of isopentenyl pyrophosphate and isoprenoids. Specifically, the current invention is directed toward identification of new methods for the synthesis of isopentenyl pyrophosphate, as isopentenyl pyrophosphate represents the universal precursor to isoprenoid synthesis."}
{"text": "This invention relates to microscopy, and particularly to scanning microscopes for confocal fluorescence microscopy.\nIn fluorescence microscopy, a sample to be examined is treated with one or more dyes that selectively attach to parts of the sample and, when illuminated with excitation light of one wavelength, fluoresce at light of another, emission wavelength. The fluorescent light is examined microscopically to identify the parts of the sample to which the dyes have attached. Typically, in modem fluorescence microscopy, the illumination is provided with a laser, which provides relatively intense light over a narrow spectrum to excite particular dyes selectively.\nIt is advantageous in fluorescence microscopy to use a scanning microscope that images only one point or pixel of the sample at a time. Where the image is to be digitized for computerized enhancement or analysis, this permits descritization of the sample by use of a relatively simple detector and relatively simple optics. Also, scanning microscopy makes confocal microscopy possible, which provides higher resolution than would otherwise be available. As discussed hereafter, scanning microscopy may be accomplished either by moving the scanning spot using a fixed objective lens, by moving the optical system itself, or by moving the sample.\nAn early approach to scanning microscopy is illustrated in Petran et al. U.S. Pat. No. 3,517,980, which discloses the use of a fixed, extended light source focused on a Nipkow disc, having spiral apertures, which produces a point source of light that scans the object field of a fixed microscope objective. The light emitted from the sample is focused on a complementary aperture in the disc and thereafter propagated to an eyepiece to effectuate confocal microscopy. A similar approach is disclosed in McCarthy et al. U.S. Pat. No. 4, 863,252, which is directed to producing relative movement of the optics and sample stage for selecting the region and depth of the sample to be scanned.\nMore recently, scanning microscopes have employed a fixed point source and fixed objective, and cause the illumination light and the light emitted from the specimen to be scanned. For example, Hayashi U.S. Pat. No. 6,028,306 discloses a scanning microscope wherein the beams produced by a plurality of point sources are caused to scan in two dimensions by pivoting mirrors operated by galvanometers. More generally, Atkinson U.S. Pat. No. 6,108,127 describes the use of scanning optics to alter the path of a light beam through a fixed microscope objective.\nThe aforementioned types of scanning microscopes require highly corrected and expensive lenses to achieve high resolution and sensitivity. They also require complex and expensive optical systems mechanical scanning devices.\nAnother recent approach to scanning microscopy is illustrated by Boettner et al. U.S. Pat. No. 5,880,465, wherein the microscope objective itself is translated laterally in two orthogonal directions to scan the sample by moving the focal point of the objective. While this approach provides the advantage of reducing the effect of off-axis light beams by ensuring that the light emitted from the specimen is always on axis, the use of linear translation to cover the entire region to be scanned limits both the area that can be scanned and the speed at which scanning can occur, and can work only with low mass objectives which have limited optical performance.\nAccordingly, there is a need for a scanning microscope that provides high resolution, high sensitivity and high scanning speed using relatively straightforward and inexpensive, but high performance, optics and mechanics.\nThe invention of the present application meets the aforementioned need by providing a scanning microscope that employs a combination of rotational and linear scanning to permit the use of a relatively simple optical elements and a relatively simple mechanical scanning system. In its general form the microscope comprises an objective lens for receiving light emitted from a sample in object space and propagating it to image space thereof, a collection lens that receives light from the objective lens and propagates it to a focal point in image space of the collection lens, and a motor, having an axis of rotation that is offset from and extending in substantially the same direction as the optical axis of the objective lens, for rotating the objective lens around the axis of rotation. This causes an optical pathway to be scanned across a sample in object space of the objective lens.\nThe sample is mounted on a stage. After each rotation of the objective lens, the stage and the optics are moved relative to one another with respect to the axis of rotation so that each subsequent scan covers a new part of the sample. Thus, the entire sample may be scanned. The stage and the optics may be moved relative to one another in another lateral direction as well, to move another sample, or sample portion, under the scanning area.\nFor fluorescence microscopy, a light source, typically a laser, is provided. A wavelength-selective beamsplitter directs the laser light toward the objective lens, while allowing fluorescence or reflected light emitted from the sample to pass through to the collection lens.\nA photo detector is placed at the focal point of the collection lens in image space so as to convert the intensity of the light emitted from the sample at the focal point of the objective lens in object space to an electrical signal. Preferably, an aperture is placed at the image space focal point of the collection lens so as to provide confocal microscopy. Additional detectors may be used, combined with corresponding wavelength-sensitive beamsplitters to tap light from the collection lens, for measuring the intensity of various wavelengths of fluorescence light produced by excitation of the sample.\nTo achieve high resolution, the maximum numerical aperture of the objective lens may be used. A light-path-altering optical element is inserted between the light source and the objective lens so as to translate the light centered on the axis of rotation to light centered on the optical axis of the objective lens, thereby filling the entrance pupil of the objective lens despite rotation thereof.\nWhere the objective lens is fixed, a rotating light-path-altering optical element receives light from an off-objective-optical-axis point in object space of said objective lens and directs it along the objective optical axis, and vice versa.\nAccordingly, it is a principal object of the present invention to provide a novel and improved scanning microscope and scanning microscopy method.\nIt is another object of the invention to provide a scanning microscope that achieves high resolution and sensitivity with relatively straightforward optics.\nIt is a further object of the invention to provide a scanning microscope that achieves high resolution and sensitivity with a relatively straightforward scanning mechanism.\nIt is yet another object of the invention to provide a relatively high speed scanning microscope.\nThe foregoing and other objects, features, and advantages of the invention will be more readily understood upon consideration of the following detailed description of the invention, taken in conjunction with the accompanying drawings."}
{"text": "1. Field of Invention\nThis invention relates to semiconductor light emitting devices including a wavelength-converting material and a filter material.\n2. Description of Related Art\nThe color of light emitted from a semiconductor light emitting device chip such as a light emitting diode may be altered by placing a wavelength-converting material in the path of the light exiting the chip. The wavelength-converting material may be, for example, a phosphor. Phosphors are luminescent materials that can absorb an excitation energy (usually radiation energy) and store this energy for a period of time. The stored energy is then emitted as radiation of a different energy than the initial excitation energy. For example, “down-conversion” refers to a situation where the emitted radiation has less quantum energy than the initial excitation radiation. The energy wavelength effectively increases, shifting the color of the light towards red.\nIf some light emitted from the chip is not absorbed by the phosphor, the unconverted light emitted from the chip mixes with the light emitted from the phosphor, producing a color between the color of the light emitted from the chip and the color of the light emitted from the phosphor. When used in applications requiring a particular color, the color of light emitted from the chip and the amount of light converted by the phosphor must be tightly controlled."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor technology. More particularly, the present invention relates to an improved method of fabricating a storage node with a supported structure of a stacked capacitor device.\n2. Description of the Prior Art\nIt has been the trend to scale down the sizes of memory cells to increase the integration level and thus memory capacity of a DRAM chip. As the sizes of DRAM devices are decreased, the capacity of the capacitors in the DRAM devices is correspondingly decreased. One approach to increasing capacity of the capacitor involves increasing the surface area of the storage node. As known in the art, the surface area of a storage node in a capacitor-over-bit-line (COB) structure is mostly increased by increasing the height since the design rule limits the horizontal dimension of the storage node. However, a higher storage node height causes instability of storage node structure, which is the cause of device failure due to two-bit or multi-bit failure during DRAM operation.\nFIGS. 1-5 are schematic, cross-sectional diagrams showing a conventional method for fabricating a storage node of a crown-type stacked cell capacitor. As shown in FIG. 1, a substrate 10 such as a silicon substrate having thereon conductive blocks 12a and 12b is provided. A dielectric layer 14 such as silicon nitride and a dielectric layer 16 such as undoped silicate glass (USG) are deposited over the substrate 10.\nAs shown in FIG. 2, a conventional lithographic process and a dry etching process are carried out to etch high aspect ratio openings 18a and 18b into the dielectric layers 14 and 16. Subsequently, a cleaning process may be performed to remove the etching byproducts or particles from the surfaces of the substrate 10 and from the interior surfaces of the openings 18a and 18b. \nAs shown in FIG. 3, a chemical vapor deposition (CVD) process is carried out to form a conformal silicon layer 22 on the surface of the dielectric layer 16 and on the interior surfaces of the openings 18a and 18b. The silicon layer 22 may be doped polysilicon.\nAs shown in FIG. 4, a planarization process such as chemical mechanical polishing (CMP) is performed to selectively remove the silicon layer 22 from the surface of the dielectric layer 16, while leaving the silicon layer 22 on the interior surfaces of the openings 18a and 18b intact.\nSubsequently, as shown in FIG. 5, a wet etching process involving the use of HF/NH4F chemistry or Buffer Oxide Etcher (BOE) is performed to remove the dielectric layer 16, thereby forming storage nodes 30a and 30b. Typically, the height H of the storage nodes 30a and 30b is approximately equal to the depth of the openings 18a and 18b, which is normally 1.6-1.7 micrometers.\nOne drawback of the above-mentioned prior art method is that when etching the high aspect ratio openings 18a and 18b it is difficult to obtain a straight sidewall profile. The tapered sidewall profile of the high aspect ratio openings 18a and 18b leads to small bottom critical dimension A. The small bottom critical dimension A results in so-called storage node bridge phenomenon during subsequent cleaning or drying processes."}
{"text": "2-Lower alkyl-cyclopentan-1,3-diones are important intermediates for the total synthesis of steroids. A number of syntheses of this important class of intermediates have been reported, including syntheses which involve the cyclization of gamma-ketoesters in the presence of base. In particular, U.S. Pat. No. 3,349,130 to Bucourt, et al., teaches such cyclization in the presence of alkali metal tertiary alcoholates. French Pat. No. 1,564,830 to Watson, et al. teaches such cyclizations in the presence of sodium hydride.\nSince alkali metal tertiary alcoholates and sodium hydride are relatively expensive and somewhat difficult to prepare in large quantities and are thus not practical for a large-scale commercial cyclization process, it would be desirable to have a cyclization method which would afford high yields of product and which would utilize relatively inexpensive and easy to prepare cyclization agents."}
{"text": "In existing communication networks, end-users may request network-based content (e.g., text, images, audio, video, and the like) available from media content servers. Disadvantageously, however, existing media content servers may be slow and, therefore, unsuitable for high-bandwidth multimedia services (e.g., online video rentals). Furthermore, since such high-bandwidth multimedia services use a combination of content streaming and storing, end-user devices with limited storage and processing capabilities (e.g., cellular phones) cannot support many high-bandwidth multimedia services (e.g., mobile music and video downloads). Moreover, existing multimedia storage solutions do not provide back-up services for storing end-user generated content."}
{"text": "1. Field of the Invention\nThis invention relates to novel fungicidal diacylimide compositions which are useful as agricultural fungicides.\n2. Description of the Prior Art\nMany compositions are known in the literature which are active as fungicides; however, the compounds of this invention show broad spectrum fungicidal activity, particularly against many agricultural pathogens."}
{"text": "Inflatable elements, bag or belt, deploying from locations adjacent vehicle occupants have been proposed and suggested to distribute belt loading during a collision (U.S. Pat. Nos. 3,682,498 and 3,841,654).\nPrior restraint systems have combined seat belts, including lap and shoulder components, with inflatable members. For example, vehicle air bags have been proposed to be mounted adjacent shoulder belts and lap belts for deployment upon rapid deceleration of a vehicle (U.S. Pat. No. 5,062,662). Other prior inflatable bag vehicle restraint systems have required that the bag be supported by a portion of the vehicle in front of the occupant (i.e., the dashboard or wheel post unit). Further, prior lap belt mounted bags were deployable in front of the occupant's belt and have not caused the lap belt to have its slack removed by the inflation of the bag.\nFinally, it has been proposed to provide bags for inflation between the occupant and shoulder straps (U.S. Pat. No. 3,971,569).\nNone of the prior art proposals provide proper protection where the restraint system can only be deployable from and restrained by a lap belt area."}
{"text": "This invention relates to water swellable absorbent articles made from crosslinked polyelectrolytes, methods for their preparation, and to a composition consisting of polyelectrolytes containing a copolymerized crosslinker which is useful to make absorbent articles.\nIt is known from U.S. Pat. Nos. 3,669,103 and 3,670,731 that cross-linked polymeric sorbents can be sandwiched between flexible supports to achieve disposable diapers or dressings.\nIt is further known from U.S. Pat. Nos. 2,988,539; 3,393,168; 3,514,419 and 3,557,067 that water swellable cross-linked carboxylic copolymers can be prepared. However, these prior art copolymers are all crosslinked during copolymerization or crosslinked after polymerization with subsequent neutralization of the carboxylic acid groups to form water swellable polyelectrolytes and hence these prior art polyelectrolytes cannot be crosslinked in situ as a coating on a substrate or as a flexible film thereof.\nIt is known from Ser. No. 468,794 filed Apr. 9, 1974 and Ser. No. 450,650 filed Mar. 13, 1974 that water swellable cured articles can be made from polyelectrolytes that have been crosslinked after polymerization with the addition of special crosslinking compounds prior to the heating or curing step.\nThe present invention is an improvement over these inventions in that the crosslinking agent is built into or copolymerized with the polyelectrolytes during polymerization. The advantage of this invention is that only a single solution needs to be sent to users and the chance of operator error in the mixing of the crosslinking agents is eliminated."}
{"text": "Generally, modern automatic battery chargers employ one of a number of charging methods, such as demand-based automatic boost charging, time-limited demand-based automatic boost charging, fixed extension demand-based automatic boost charging, or the like.\nLimitations of standard demand-based boost charging systems and methods include that initiation and reversion thresholds must exceed the DC load on the battery plus the fully charged battery's current at boost charging voltage. Another limitation is that the initiation and reversion thresholds must allow for battery current variations caused by battery manufacturing differences, temperature, and aging. Because the method is not adaptive, another limitation of standard systems is that thresholds must be selected to avoid over-charging for the severe cases. This results in under-charging in most applications.\nLimitations of time-limited demand-based boost charging systems and methods include the same general limitations as the standard demand-based boost charging method. The time-limited systems adds a backup means of boost charge termination, limiting the impact of over-charging if the initiation and reversion thresholds are too optimistic for the application. Another limitation is the time limit is not adaptive. It is set to a time that is ideal for only one specific recharge case and non-optimal for all others.\nLimitations of fixed-time-extension demand-based boost charging systems and methods have the same general limitations as the standard demand-based boost charging method. It adds a fixed time extension to the boost charge, for the purpose of achieving a more complete charge for conservative reversion transition current settings. The limitation is that time extension is not adaptive. The time extension is the same regardless of battery condition, depth of battery discharge, temperature, etc. The selected value is a compromise between over-charging in some cases and under-charging in others. For example, some existing battery charges have a fixed six-hour extension to a demand-based boost charge. Since this time extension is fixed the battery will be over charged when it is lightly discharged and requires a shorter boost charge, or it will be under-charged when it is deeply discharged and requires a longer boost charge."}
{"text": "1. Field of the Invention\nThis invention relates to a multi-articulated arm mechanism suitable for use in working robots and manipulators, and more particularly to an arm mechanism which has high rigidity against operational reaction forces or other loads which are imposed on the arm during operation.\n2. Description of the Prior Art\nIn general, a multi-joint robot has a series of arms with articulated joints on a fixed base or pedestal, performing an assigned job by a working tool which is attached to the distal end of the arms. Therefore, the force applied to a work is invariably reflected by a bending moment acting on the base portion of the arms, causing operational errors by arm flexure which occurs to a greater degree with longer arms. For this reason, the application of the multi-articulated robots has been limited to those fields in which the robots can perform a required job without contacting the work, free of the reaction forces as in welding and coating operations. They have been considered unsuitable for operations like grinding or other machining operations which involve reaction forces or operations which require attaching a testing device or other devices of a large weight at the distal end of the arm.\nNevertheless, the multi-articulated arm mechanism is superior in operability to the arm mechanisms of the perpendicular coordinates system or the parallel type arm mechanisms, and accordingly extremely convenient for performing complicated operations, for example, the operations which have to be carried out in a narrow limited space or the operations which require evasion of obstacles."}
{"text": "The present disclosure relates to a light emitting device, a method of manufacturing the same, a light emitting device package, and a lighting system.\nA Light Emitting Diode (LED) is manufactured by combining a p-n junction diode, which is characterized in converting electric energy into light energy, with Group III and V elements in the periodic table of the chemical elements. The LED may produce various colors by adjusting a composition ratio of a compound semiconductor and materials.\nThe LED may generate light energy corresponding to an energy gap between a conduction band and a valance band by combining electrons in an n-layer and electron-holes in a p-layer when forward voltage is applied.\nEspecially, a blue LED, a green LED, and an Ultra Violet (UV) LED using a nitride semiconductor are commercialized and are extensively used."}
{"text": "The present invention relates to a pull-on disposable wearing article.\nJapanese Unexamined Patent Application Publication No. 1993-317356 (REFERENCE) discloses a pull-on/closed type diaper comprises a pair of fastening strips bonded to transversely opposite side edges of front and rear waist regions so as to extend outward and tear-apart lines provided immediately inside the bonded region in each of the fastening strips along which the front and rear waist regions are disconnected from each other. A distal end of the fastening strip is provided with an engaging region defined by a pressure-sensitive adhesive coated on this distal end or defined by a tape strip provided with a plurality of hooks bonded thereto. This disposable diaper provided with the tear-apart lines and the fastening strips allows the diaper to be used as an open-type diaper after the front and rear waist regions have been disconnected from each other along the tear-apart lines.\nWhile the fastening strips are not actually used before the front and rear waist regions are disconnected from each other along the respective tear-apart lines, there is anxiety that the engaging region might be mischievously contaminated to unworkable condition and/or might come in contact with the wearer's skin and might uncomfortably irritate the wearer's skin."}
{"text": "1. Field of the Invention\nThe present invention relates to an asynchronous serial data receiving device and an asynchronous serial data transmitting device, which are connected to or installed in a data processing device such as a microcomputer or the like, and more particularly relates to an asynchronous serial data receiving device and an asynchronous serial data transmitting device, which operate data communication based on the asynchronous data receiving and transmitting method capable of controlling the data reception conducted at the receiving device in the case where sampling interval between each of the bits forming a communication data is not constant, and also controlling the data transmission conducted at the transmitting device in the case where transmission interval between each of the data forming bits is not constant.\n2. Description of the Related Art\nMethods of communication between different microcomputers or between a microcomputer and a peripheral device are classified into two methods; namely a method in which data are received and/or transmitted in parallel per each data unit formed by a plurality of bits, and a method in which data are received and/or transmitted in serial per each data bit. The method in which data are received and/or transmitted in parallel can communicate more data per unit time than the case in which data are received and/or transmitted in serial. However, the former method requires many more cable distributions, which thereby increases the total cost for communications overall.\nOn the other hand, the methods in which data are received and/or transmitted in serial can be further classified into two communication methods; namely the synchronous serial data transmission and/or reception method in which the timing for synchronization between the transmission side and the reception side is adjusted by use of a clock signal, and the asynchronous serial data transmission and/or reception method in which no clock signal is used for adjusting the timing for synchronization. Since the asynchronous serial data transmission and/or reception method requires no clock signal, although it requires less cable distributions than those of the synchronous serial data transmission and/or reception method, its operation speed tends to be slow.\nFIG. 23 is an exemplary view showing the general idea of a receiver section and a transmitter section in the asynchronous serial data communication system. In the figure, reference numeral 220 denotes an asynchronous serial data transmitting device as the transmitter section in the asynchronous serial data communications, numeral 221 denotes an asynchronous serial data receiving device as the receiver section in the asynchronous serial data communications, numeral 223 denotes a communication enable signal which is output from the receiver section 221, and input to the transmitter section 220, and reference numeral 224 denotes a communication data which is output from the transmitter section 220, and input to the receiver section 221.\nThe operation in the conventional asynchronous serial data communication system is now explained below.\nFIG. 24 is an exemplary view showing an example of the communication data format used in the asynchronous serial data communications. In the figure, reference numeral 5 denotes a start bit which is a one-bit signal of the logic low level, 6 denotes a data bit, 7 denotes a parity bit which is added to the communication data in order to improve reliability of the data, and reference numeral 8 denotes a stop bit for indicating the end of the data transmission, which is formed by either one-bit signal or two-bit signal of the logic high level.\nxe2x80x9cParityxe2x80x9d includes an even parity which sets the parity bit in such a manner that the number of bits of the logic xe2x80x9c1xe2x80x9d (hereinafter may be referred to as xe2x80x9clogic high levelxe2x80x9d or just as xe2x80x9chighxe2x80x9d) in the communication data formed by the total number of data bits 6 and parity bit 7 becomes a certain even number, and also includes an odd parity which sets the parity bit in such a manner that the number of bits of the logic xe2x80x9c1xe2x80x9d in the communication data formed by the total number of data bits 6 and parity bit 7 becomes a certain odd number. There is also communication data which requires no parity bit.\nFIG. 25 is a block diagram of the conventional asynchronous serial data receiving device as the receiver section 221. In the figure, reference numeral 252 denotes a serial-to-parallel conversion circuit which inputs a communication data 224 on the basis of a data shift signal 254, and converts it from the serial data to the parallel data 253.\nFIG. 26 is a timing chart indicating the operation of the serial-to-parallel circuit 252 shown in FIG. 25. When the receiver section 221 is in the state in which it can receive the data signal, it sends a communication enable signal 223 by setting it to the logic low level to the transmitter section 220. The transmitter section 220 recognizes that the communication enable signal is in the low level, and sends a binary data signal xe2x80x9c0101001001xe2x80x9d.\nOn this occasion, the start bit 5 has been added to the head portion of the binary data signal to be transmitted, and also the stop bit 8 has been added to the last portion of the data. The receiver section 221 receives the communication data 224, detects the start bit 5 and starts receiving the entire communication data 224.\nThe same baud rate is set to both the transmitter section 220 and the receiver section 221, wherein the reception and transmission of the communication data are executed in accordance with the thus set baud rate.\nUnder the ideal condition, the value of the communication data 224 is taken into the serial-to-parallel conversion circuit 252 at the center of each bit. In FIG. 26, the communication data 224 is taken into the serial-to-parallel conversion circuit 252, when the data shift signal 254 is changed from high to low.\nSince the conventional asynchronous serial data transmitting and receiving devices are configured as such, the transmitter section 220 and the receiver section 221 generally operate at different operation clocks from each other. For this reason, since the time required for communicating a data of one bit can be set only to a time period whose length is a multiple of integer of an operation clock cycle, even in a case where the same baud rate is set to both the transmitter section 220 and the receiver section 221, there has been caused a setting error of the baud rate.\nFIG. 26 shows an example of the operation in the case where the receiver section 221 receives a communication data formed by eight data bits 6 and one data bit 8, when the communication time for one-bit data is set to the time period corresponding to four cycles of the clock signal 11 used therein.\nIn the figure, it is arranged in such a manner that when the data shift signal 254 is changed from the high level to the low level, the communication data 224 is taken into the serial-to-parallel conversion circuit 252, and the data taken into the serial-to-parallel conversion circuit 252 is formed as xe2x80x9c101000100xe2x80x9d excluding the start bit 5.\nIn the example shown in FIG. 26, a communication data receiving error can be observed at the sixth, seventh and ninth bits of the communication data 224. If the baud rate set in the case of FIG. 26 is made smaller, the communication time for one-bit data is made longer, so that occurrence of the baud rate setting error with respect to the communication time for one-bit data can be efficiently suppressed.\nFor this reason above, in the actual case of the asynchronous serial data transmitting and receiving devices, there has been such a problem that the largest baud rate has to be suppressed down to the level in which no data receiving error occurs.\nThe present invention has bee n proposed to solve the problems aforementioned, and it is an object of the present invention to provide an asynchronous serial data receiving device and an asynchronous serial data transmitting device, which operate data communication based on the asynchronous data receiving and transmitting method capable of raising the maximum baud rate, yet even without changing the frequency of the operation clocks at the respective receiver section and the transmitter section, by enabling the setting of baud rate of the asynchronous serial data receiving device as the receiver section and that of the asynchronous serial data transmitting device as the transmitter section more precisely.\nIn order to achieve the above object, an asynchronous serial data receiving device according to the first aspect of the present invention comprises:a serial-to-parallel data conversion means, which takes in a communication data of one transmission unit starting with a start bit, and converts and outputs the communication data into a parallel data, a control signal generating means for generating a data shift signal to indicate the timing at which the serial-to-parallel conversion means should take in each of the bit data forming the communication data, and outputting the generated data shift signal to the serial-to-parallel conversion means, a count correction means, which generates a count value correction signal for controlling the time interval between a first timing at which it takes the first bit data of the received communication data into the serial-to-parallel conversion means and the second timing at which it takes in the second bit data following the first bit data into the serial-to-parallel conversion means for each of the bit data of the communication data, and outputs the generated count value correction signal to the control signal generating means, wherein the control signal generating means controls the output timing of the data shift signal in accordance with the count value correction signal output from the count correction means, and the serial-to-parallel conversion means takes in each of the bit data of the communication data on the basis of the data shift signal.\nThe count correction means in the asynchronous serial data receiving device as constructed above comprises: a counting means for sequentially inputting the data shift signals and counting the number of the data shift signals, and a selecting means for selecting either one of the signal output from the counting means and a predetermined value, and outputting the selected value to the control signal generating means as a count value correction signal.\nAn asynchronous serial data transmitting device according to the second aspect of the present invention comprises: a parallel-to-serial data conversion means, which transmits a communication data of one transmission unit starting with a start bit, and converts and outputs the communication data into a serial data, a control signal generating means for generating a data shift signal to indicate the timing at which the parallel-to-serial conversion means should transmit each of the bit data forming the communication data, and outputting the generated data shift signal to the parallel-to-serial conversion means, a count correction means, which generates a count value correction signal for controlling the time interval between a first timing at which it transmits the first bit data of the received communication data from the parallel-to-serial conversion means and the second timing at which it transmits the second bit data following the first bit data from the parallel-to-serial conversion means for each bit data of the communication data, and outputs the generated count value correction signal to the control signal generating means, wherein, the control signal generating means controls the output timing of the data shift signal in accordance with the count value correction signal output from the count correction means, and the parallel-to-serial conversion means transmits each of the bit data of the communication data on the basis of the data shift signal.\nThe count correction means in the asynchronous serial data receiving device as constructed above comprises: a counting means for sequentially inputting the data shift signals and counting the number of the data shift signals, and a selecting means for selecting either one of the signal output from the counting means and a predetermined value, and outputting the selected value to the control signal generating means as a count value correction signal."}
{"text": "Various known coffee making appliances involve adding of coffee beans in one of various forms (i.e., ground or unground) to a container that is part of a machine in which heated water is delivered to the container and passes therethrough. The container typically includes a filtering mechanism so that heated water exiting the container is in the form of brewed coffee. Certain appliances require pre-ground coffee beans to be added in the form of “grounds.” Other appliances are designed to accept whole coffee beans into a hopper or opening and include mechanisms that grind the beans into a ground form and then complete the brewing process. It is sometimes preferred to keep whole beans on hand for making coffee so that the coffee may be ground immediately prior to brewing, as this is believed to produce the freshest tasting coffee beverage. Yet other appliances required use of pods or pre-packaged coffee grounds in a serving-sized filter container to be inserted into the appliance. In such instances, the pre-packaged pods or containers are limited to being filled with ground coffee. Certain pre-packaged containers require placement into a separate filter mechanism in the appliance, while others have built-in filtering mechanisms."}
{"text": "One type of data store that may conventionally be used to implement storage and maintenance of large datasets is a block-based data store. In a block-based data store, data such as key-value records are stored in data blocks that are optimized for sequential access from storage. In order to enable the performance of queries of a block-based data store, indices may be constructed that are associated with various attributes, values, or other information in the data store. To date, the development of indexing techniques for block-based data stores has focused mainly on maximizing query speed for complete, static sets of data blocks. Thus, conventional techniques commonly regard indexing as a one-time process performed upon importation of a complete dataset. Such conventional techniques typically do not prioritize indexing speed and/or efficiency, and may not enable efficient incremental indexing of dynamic datasets.\nIn a wide variety of real-world applications, indexing speed and/or efficiency may be important considerations, as may be the ability to incrementally index a dataset that is constantly changing. Some dynamic block-based data stores may continually receive and incorporate new records, in some cases millions or even billions of times per second. As such, techniques for efficient incremental indexing of block-based data stores may be desirable."}
{"text": "Platinum is a candidate for utilization as a conductive material in advanced semiconductor processing. Platinum can be utilized in an elemental form, or as an alloy (such as, for example, rhodium/platinum), and can be deposited onto a substrate by, for example, sputter deposition or chemical vapor deposition (CVD) methods. Platinum is typically formed to have a relatively smooth upper surface. Such smooth upper surface can be advantageous in, for example, applications in which circuitry is formed over the platinum layer. Specifically, the relatively smooth surface can provide a substantially planar platform upon which other circuitry is formed. However, there can be advantages to incorporating roughened conductive layers into integrated circuitry in applications where high surface area is desired, as with capacitor electrodes. Accordingly, it would be desirable to develop methods of forming platinum layers having roughened outer surfaces.\nIn another aspect of the prior art, platinum-comprising materials are frequently utilized as catalysts in, for example, the petroleum industry, as well as in, for example, automobile exhaust systems. Frequently, an efficiency of a catalyst can be improved by enhancing a surface area of the catalyst. Accordingly, it would be desirable to develop methods of enhancing surface area of platinum-comprising materials."}
{"text": "The present invention relates generally to a method of using standard .ZIP files and strong encryption technology to securely store files, and more particularly to a method of integrating existing strong encryption methods into the processing of .ZIP files to provide a highly secure data container which provides flexibility in the use of symmetric and asymmetric encryption technology. The present invention adapts the well established and widely used .ZIP file format to support higher levels of security and multiple methods of data encryption and key management, thereby producing an efficient, highly secure and flexible digital container for electronically storing and transferring confidential data.\nCompression of computer files has been available for many years. Compressing files can save large amounts of disk space, and can reduce transfer time when downloading files from the Internet or transferring files through email. Almost any file one downloads from the Internet is compressed in some way. A standard compressed file or folder as it is sometimes called contains one or more files that were compressed into a single file or folder. Many different compression formats have been developed over the years. The .ZIP format, created by the assignee of the present invention, is perhaps the most common compressed file format for the personal computer. Any file with a “.zip” extension most likely contains one or more files of data archived, that is, each either compressed or stored, in the .ZIP format. “.Zipping” a file has become a commonly used term meaning to compress the file into the .ZIP format archive so that it occupies less disk space, and similarly, “unzipping” a file means decompressing a compressed file in the .ZIP format.\nA .ZIP file is generally recognized as a data compression and archiving format invented by PKWARE, Inc. The .ZIP format is a file format designed for combining data compression technology with file archiving techniques. Many commercially available software products are available for compressing or “zipping” files or other data into the .ZIP format. These .ZIP files can then be used to reconstruct the original data through the “unzipping” process. Data compression converts the contents of a file into an encoded format requiring less computer storage space or in the case of transmission less network bandwidth than the original uncompressed file.\nArchiving, in the context of a .ZIP file, is a method of storing information about the characteristics of a file in a catalogue of files, known as the Central Directory, inside the .ZIP file, allowing each file to be retrieved individually by its characteristics. This capability is widely used. These characteristics include, but are not limited to, file name, file size, and file creation date and time.\nSoftware programs such as PKZIP® written by PKWARE, Inc. are used to process files in the .ZIP format. Such programs allow one or more files of any type to be compressed and archived into a file of the .ZIP format type for efficient file storage and transmission over computer and communication networks. This format and the software programs that process .ZIP files have become ubiquitous.\nData encryption is used by many software programs to provide data privacy. Data encryption is a method of encoding data so that it cannot be reproduced in its original form unless an associated key is provided. Decryption uses this key to convert the encrypted data back into its original state. The key is known only to the person encrypting the data or by those other people with whom the person encrypting the data chooses to share the key. The key is used to “unlock” the data so that it can again be used in its original form.\nKeys are uniquely generated using data known to the person encrypting a file or other data associated with recipients and users of the file. This data can be a user-defined password or other random data. Several methods are commonly used for processing the keys used for data encryption. Encryption using a key generated from a password is an example of symmetric encryption. Encryption using a public/private key pair is an example of asymmetric encryption. An example of one method for processing encryption keys supported by this invention uses a public/private key pair commonly associated with digital certificates as defined by the document Internet X.509 Public Key Infrastructure Certificate and CRL Profile (RFC 2459). A digital certificate is a unique digital identifier associating a public and private key pair to an assigned individual, a group, or an organization. When used for encrypting data, the public key of an individual is used to process an encryption key which only the individual in possession of the corresponding private key can use for decryption. A digital certificate is issued to an individual, a group, or an organization for a fixed period of time and can only be used during this time period. After the time period has elapsed, the digital certificate will be considered to have expired and must be reissued for a new time period.\nThe strength of a data encryption method is determined at least in part by its key size in bits. The larger the key size a data encryption method uses, the more resistant it is to cryptanalysis. Cryptanalysis, or popularly “cracking”, is the unauthorized access to encrypted data. Strong encryption is a type of data encryption that uses key sizes of 128 bits or more. A number of encryption encoding methods are known today. Examples supported by the present invention include but are not limited to Advanced Encryption Standard (AES), Data Encryption Standard (DES), 2DES, 3DES, and others. A number of key sizes are commonly used today. Examples supported by the present invention include but are not limited to 128 bits, 192 bits, and 256 bits.\nMany software programs available today that process .ZIP files use data encryption to encrypt files after compression as they are written to the .ZIP file. The data encryption method used by these software programs uses a key size of 96 bits or less and is considered weak or moderate encryption by today's standards. These software programs use keys generated using user-defined password data. Weak data encryption may not provide sufficient security to computer users that store and transfer their confidential data files using the .ZIP format.\nPassword-based key generation has been a commonly used method of applying data encryption, however, known vulnerabilities to cracking methods such as “brute force password cracking” make this method of encryption insufficient to meet today's more advanced security needs. Another known limitation of password-based security is the lack of non-repudiation. Non-repudiation is the ability to be certain that the person or program that created an encrypted .ZIP file cannot deny that fact and that their identity is bound to the .ZIP file they created. This cannot be achieved with symmetric encryption methods. Today, non-repudiation is an important aspect of security related to the implementation of digital certificates and digital signatures. It is critically important to be able to prove that a creator or sender of an encrypted file did in fact create the file, i.e. not repudiate his/her action.\nTherefore, a need exists to extend the options for levels of security available to programs that process .ZIP files. This extended of security capability makes use of the encryption technologies available today or others that may gain acceptance in the future."}
{"text": "1. Field\nApparatuses consistent with example embodiments relate to a plasma monitoring apparatus and a plasma processing apparatus including the same.\n2. Description Of Related Art\nIn the related art, a semiconductor device is manufactured through a plurality of unit processes including, for example, a deposition process using plasma, an etching process, and the like. In accordance with the miniaturization of semiconductor devices and high degree of integration of semiconductor products, the effects of a non-uniform distribution of plasma to the properties of semiconductor devices are increased."}
{"text": "This invention relates generally to coin handling and processing apparatus such as coin wrapping apparatus and more particularly to a coin stacking tube device suitable for incorporation in any of such apparatus.\nA typical example of a conventional coin wrapping apparatus is that of a construction wherein coins supplied from a hopper onto a turntable are arrayed by centrifugal force applied thereto along the periphery of the turntable, and a coin counting mechanism counts the coins while the coins are sent from the turntable to a coin passage and propelled therealong by, for instance, a propelling belt. The coins delivered from the coin passage one after another are received in a coin stacking tube to be stacked therein, the coins thus stacked being dropped into a coin wrapping device comprising a plurality of coin wrapping rolls disposed around a circle directly below the coin stacking tube. The thus stacked coins are wrapped by the wrapping rolls with a piece of paper, and the lateral edges of the paper are fold crimped to form firm beads by which the paper is maintained in tightly wrapped state.\nIn the case when it is required to change the denomination of coins to be wrapped in the coin wrapping apparatus, various parts of the apparatus must be readjusted so that the apparatus is set for the new denomination of coins. Of these parts, the number of coins supplied onto the turntable, the rotating speed of the same, the lateral width of the coin passage, the height and the position of the propelling belt, the positions of the wrapping rolls to be brought into contact with the stack of coins, and the rotating speed of the wrapping rolls can be readjusted comparatively easily by interlinking those members controlling these values with a member for setting the apparatus to a different denominations of coins.\nHowever, the adjustment of the inner diameter of the coin stacking tube is not easy, and therefore it has been a conventional practice to prepare a number of coin stacking tubes each having a different inner diameter suitable for a specific denomination (or one outer diameter) of coins and, at the time of changing the denomination of coins to be wrapped, to replace the existing coin stacking tube with another suitable for the new denomination.\nIn this case also, the denomination changing operation is found to be troublesome because of the requirement of the selection and the replacement, and furthermore there has been a high possibility of erroneous selection of the coin stacking tube, causing unsatisfactory stacking of coins in the erroneous coin stacking tube.\nRecently, a type of coin stacking tube whose inner diameter is made variable in accordance with the denomination of coins has been developed. In this kind of coin stacking tube, the space for receiving the coins is formed by a plurality of blades provided in an outer casing, the inner edge of each blace contacting against the inner surface of the preceding blade, the stem part of each blade being supported rotatably about an axis fixed to the outer casing, and all blades being rotated in either of the opening and closing directions in accordance with the denomination setting in the coin wrapping apparatus, whereby the space formed within the plurality of blades is adapted for the denomination or the outer diameter of the coins to be stacked therein.\nIn this example, however, in order to assure smooth rotation of the blades at the time of expension or constraction of the interior space of the coin stacking tube, the rotating axes for the blades must be secured to the outer casing in parallel with each other. If the rotating axes are not accurately, parallel, the coin stacking space formed therein cannot be of a correct circular cross-section, thus causing unsatisfactory stacking of coins within the tube, and the uneven contacting of the innermost edges of the blades against the inner surfaces of the preceding blades causing a considerable torque to be required for rotating the blades."}
{"text": "1. Field\nMethods and apparatuses consistent with exemplary embodiments relate to a 3-dimensional (3D) display device and a method thereof, and more particularly, to a 3D display device which uses a barrier to enable a user to view a 3D image without 3D glasses, and a method thereof.\n2. Description of the Related Art\nWith the development of electronic technologies, various kinds of electronic apparatuses have been developed and distributed. Display apparatuses such as televisions (TVs) are representative of home appliances.\nIn recent years, 3D display devices, which enable a user to view a 3D image, are increasingly being distributed. The 3D display devices may be divided into a glasses type system and a non-glasses type system according to whether the 3D display device uses 3D glasses.\nOne example of the glasses type system is a display device which employs a shutter glasses method. In the shutter glasses method, the display device outputs a left-eye image and a right-eye image alternately, and a left shutter glass and a right shutter glass of a 3D glasses apparatus worn by a user are opened and closed alternately in association with the output of the display device, thereby making the user feel a sense of depth.\nThe non-glasses type system may be referred to as an autostereoscopy system. A 3D display device of the non-glasses type method displays a space-shifted multi-view image, and projects light corresponding to images of different points of view onto a left eye and a right eye of a user using a parallax barrier or a lenticular lens, thereby making the user feel a sense of depth.\nAs described above, the non-glasses type system has the advantage of enabling users to view a 3D image without glasses.\nFIG. 1 is a view to explain an operation of a related-art display device which uses a parallax barrier.\nReferring to FIG. 1, a barrier 10 is disposed on an upper portion of a display panel 20. The barrier 10 includes a plurality of vertical line patterns. An odd line (a) and an even line (b) are turned on or off by being driven alternately.\nThe display panel 20 displays a frame in which a left-eye image (L) and a right-eye image (R) are placed alternately in a vertical column direction, and then displays a frame in which a left-eye image (L) and a right-eye image (R) are placed in reverse order of the previous frame.\nThe barrier 10 switches driving of the odd line and the even line according to the operation of the display panel 20. Accordingly, the left-eye image is continuously input to the left eye and the right-eye image is continuously input to the right eye, such that the user can feel a sense of depth.\nHowever, since such a non-glasses type system provides a multi-view image, crosstalk in which a left-eye image and a right-eye image are mixed may occur.\nIn particular, when the user changes his or her location, the right-eye image may be projected onto the user's left eye and the left-eye image may be projected onto the user's right eye according to the movement of the user, and thus, a reverse image in which the left-eye image and the right-eye image are reversed may be generated.\nTherefore, there is a demand for a method for displaying a 3D image effectively in a non-glasses type system."}
{"text": "I. Field\nThe present disclosure relates generally to communication, and more specifically to techniques for high rate data transmission.\nII. Background\nWireless communication systems are widely deployed to provide various communication services such as voice, packet data, broadcast, messaging, and so on. These systems may be multiple-access systems capable of supporting communication for multiple users by sharing the available system resources. Examples of such multiple-access systems include Code Division Multiple Access (CDMA) systems, Time Division Multiple Access (TDMA) systems, Frequency Division Multiple Access (FDMA) systems, and Orthogonal Frequency Division Multiple Access (OFDMA) systems.\nData usage for wireless communication systems continually grows due to increasing number of users as well as emergence of new applications with higher data requirements. However, a given system typically has limited transmission capacity, which is determined by the design of the system. A substantial increase in transmission capacity is often realized by deploying a new generation or a new design of a system. For example, the transition from second generation (2G) to third generation (3G) in cellular systems provides substantial improvements in data rate and features. However, new system deployment is capital intensive and often complicated.\nThere is therefore a need in the art for techniques to improve transmission capacity of a wireless communication system in an efficient and cost effective manner."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a solenoid valve. More specifically, the present invention relates to a method and apparatus for controlling the solenoid valve to achieve a desired flow rate.\n2. Description of the Related Art\nSolenoid valves are used in a multitude of various operations. One use is the automotive industry wherein a solenoid valve is used in conjunction with a vapor canister in a vehicle emission system. For instance, under normal operating conditions, fuel vapors from the vehicle's fuel tank are stored in the vapor canister. The canister is purged by drawing fresh air through the canister into the intake manifold of the engine. Purging the canister disrupts optimum air-fuel ratio and may result in inefficient operation of the engine. Thus, a solenoid valve is used to control the flow rate of fuel vapor being drawn from the vapor canister into the intake manifold.\nModern engines are tightly tuned for optimum operating performance. The amount of canister purge vapor entering the intake manifold is controlled by the solenoid. The solenoid valve is turned on and off or cycled based on various operating parameters. The duty cycle or percentage of time that the solenoid valve is open regulates the flow of fuel vapors being purged from the canister.\nVarious types of control systems to control or regulate the desired flow of fuel vapor from the vapor canister are known. One type of system controls operation of the solenoid valve through a duty cycle pulse width modulation. Duty cycle pulse width modulation control systems use from 5 to 100 percent of the duty cycle to vary the flow. Such systems fail to provide the flexibility necessary to control and, more specifically, regulate flow at both the low and high ends of the duty cycle. For instance, the slope of the flow rate versus percent duty cycle for optimum low end control may not be suitable for high end flow and vice versa.\nThus, it is advisable to have a control system which provides optimum flow control characteristics throughout the entire range of the duty cycle. For instance, the flow rate may be modified independently of the duty cycle wherein the overall slope of the flow rate with respect to the duty cycle may change for various percentages of duty cycle.\nAnother type of control system uses a current signal to actuate an armature so that flow is proportional to current. These control systems tend to have more hysteresis and are not useful for the initial 25 percent of the full scale signal."}
{"text": "The present invention is directed to a fan array fan section utilized in an air-handling system.\nAir-handling systems (also referred to as an air handler) have traditionally been used to condition buildings or rooms (hereinafter referred to as “structures”). An air-handling system is defined as a system that includes components designed to work together in order to condition air as part of the primary system for ventilation of structures. The air-handling system may contain components such as cooling coils, heating coils, filters, humidifiers, fans, sound attenuators, controls, and other devices functioning to meet the needs of the structures. The air-handling system may be manufactured in a factory and brought to the structure to be installed or it may be built on site using the necessary devices to meet the functioning needs of the structure. The air-handling compartment 102 of the air-handling system includes the inlet plenum 112 prior to the fan inlet cone 104 and the discharge plenum 110. Within the air-handling compartment 102 is situated the fan unit 100 (shown in FIGS. 1 and 2 as an inlet cone 104, a fan 106, and a motor 108), fan frame, and any appurtenance associated with the function of the fan (e.g. dampers, controls, settling means, and associated cabinetry). Within the fan 106 is a fan wheel (not shown) having at least one blade. The fan wheel has a fan wheel diameter that is measured from one side of the outer periphery of the fan wheel to the opposite side of the outer periphery of the fan wheel. The dimensions of the handling compartment 102 such as height, width, and airway length are determined by consulting fan manufacturers data for the type of fan selected.\nFIG. 1 shows an exemplary prior art air-handling system having a single fan unit 100 housed in an air-handling compartment 102. For exemplary purposes, the fan unit 100 is shown having an inlet cone 104, a fan 106, and a motor 108. Larger structures, structures requiring greater air volume, or structures requiring higher or lower temperatures have generally needed a larger fan unit 100 and a generally correspondingly larger air-handling compartment 102.\nAs shown in FIG. 1, an air-handling compartment 102 is substantially divided into a discharge plenum 110 and an inlet plenum 112. The combined discharge plenum 110 and the inlet plenum 112 can be referred to as the airway path 120. The fan unit 100 may be situated in the discharge plenum 110 as shown), the inlet plenum 112, or partially within the inlet plenum 112 and partially within the discharge plenum 110. The portion of the airway path 120 in which the fan unit 100 is positioned may be generically referred to as the “fan section” (indicated by reference numeral 114). The size of the inlet cone 104, the size of the fan 106, the size the motor 108, and the size of the fan frame (not shown) at least partially determine the length of the airway path 120. Filter banks 122 and/or cooling coils (not shown) may be added to the system either upstream or downstream of the fan units 100.\nFor example, a first exemplary structure requiring 50,000 cubic feet per minute of air flow at six (6) inches water gage pressure would generally require a prior art air-handling compartment 102 large enough to house a 55 inch impeller, a 100 horsepower motor, and supporting framework. The prior art air-handling compartment 102, in turn would be approximately 92 inches high by 114 to 147 inches wide and 106 to 112 inches long. The minimum length of the air-handling compartment 102 and/or airway path 120 would be dictated by published manufacturers data for a given fan type, motor size, and application. Prior art cabinet sizing guides show exemplary rules for configuring an air-handling compartment 102. These rules are based on optimization, regulations, and experimentation.\nFor example, a second exemplary structure includes a recirculation air handler used in semiconductor and pharmaceutical clean rooms requiring 26,000 cubic feet per minute at two (2) inches water gage pressure. This structure would generally require a prior art air-handling system with a air-handling compartment 102 large enough to house a 44 inch impeller, a 25 horsepower motor, and supporting framework. The prior art air-handling compartment 102, in turn would be approximately 78 inches high by 99 inches wide and 94 to 100 inches long. The minimum length of the air-handling compartment 102 and/or airway path 120 would be dictated by published manufacturers data for a given fan type, motor size and application. Prior art cabinet sizing guides show exemplary rules for configuring an air-handling compartment 102. These rules are based on optimization, regulations, and experimentation.\nThese prior art air-handling systems have many problems including the following exemplary problems:                Because real estate (e.g. structure space) is extremely expensive, the larger size of the air-handling compartment 102 is extremely undesirable.        The single fan units 100 are expensive to produce and are generally custom produced for each job.        Single fan units 100 are expensive to operate.        Single fan units 100 are inefficient in that they only have optimal or peak efficiency over a small portion of their operating range.        If a single fan unit 100 breaks down, there is no air conditioning at all.        The low frequency sound of the large fan unit 100 is hard to attenuate.        The high mass and turbulence of the large fan unit 100 can cause undesirable vibration.        \nHeight restrictions have necessitated the use of air-handling systems built with two fan units 100 arranged horizontally adjacent to each other. It should be noted, however, that a good engineering practice is to design air handler cabinets and discharge plenums 110 to be symmetrical to facilitate more uniform air flow across the width and height of the cabinet. Twin fan units 100 have been utilized where there is a height restriction and the unit is designed with a high aspect ratio to accommodate the desired flow rate. As shown in the Greenheck “Installation Operating and Maintenance Manual,” if side-by-side installation was contemplated, there were specific instructions to arrange the fans such that there was at least one fan wheel diameter spacing between the fan wheels and at least one-half a fan wheel diameter between the fan and the walls or ceilings. The Greenheck reference even specifically states that arrangements “with less spacing will experience performance losses.” Normally, the air-handling system and air-handling compartment 102 are designed for a uniform velocity gradient of 500 feet per minute velocity in the direction of air flow. The two fan unit 100 air-handling systems, however, still substantially suffered from the problems of the single unit embodiments. There was no recognition of advantages by increasing the number of fan units 100 from one to two. Further, the two fan unit 100 section exhibits a non-uniform velocity gradient in the region following the fan unit 100 that creates uneven air flow across filters, coils, and sound attenuators.\nIt should be noted that electrical devices have taken advantage of multiple fan cooling systems. For example, U.S. Pat. No. 6,414,845 to Bonet uses a multiple-fan modular cooling component for installation in multiple component-bay electronic devices. Although some of the advantages realized in the Bonet system would be realized in the present system, there are significant differences. For example, the Bonet system is designed to facilitate electronic component cooling by directing the output from each fan to a specific device or area. The Bonet system would not work to direct air flow to all devices in the direction of general air flow. Other patents such as U.S. Pat. No. 4,767,262 to Simon and U.S. Pat. No. 6,388,880 to El-Ghobashy et al. teach fan arrays for use with electronics.\nEven in the computer and machine industries, however, operating fans in parallel is taught against as not providing the desired results except in low system resistance situations where fans operate in near free delivery. For example, Sunon Group has a web page in which they show two axial fans operating in parallel, but specifically state that if “the parallel fans are applied to the higher system resistance that [an] enclosure has, . . . less increase in flow results with parallel fan operation.” Similar examples of teaching against using fans in parallel are found in an article accessible from HighBeam Research's library (http://stati.highbeam.com) and an article by Ian McLeod accessible at (http://www.papstplc.com)."}
{"text": "Adjustable water ski bindings are utilized to attach water skis to the feet of skiers and consist of a foot piece which is stationary on the water ski and a heel piece which is adjustable lengthwise of the water ski. In order to attach the water ski to the foot, the toe is first inserted into the front vamp or foot piece and then the heel piece is pressed against the heel until a tight fit on the foot is accomplished. The heel piece can be attached to a mounting plate which is slidable in guides attached to the ski on opposite sides of the plate. The mounting plate can carry two latch members having teeth which cooperate with ratchet teeth on the guides and when the teeth are engaged, the heel piece will be locked against rearward movement which would loosen the heel piece. Quick loosening of the heel piece can be accomplished by disengagement of the teeth by movement of the latch members so that the skier is able to quickly remove the ski from his foot in the event of a fall into the water. Present locking mechanisms for the heel piece utilize pawl or latch members and springs which are unnecessarily complicated since either pivots or guide slots are required in the mounting plate for the latch members and separate springs are required for these members. The U.S. Pat. to W. J. Meucci No. 3,137,014 is an example of guide slots cut in the mounting plate for latch or pawl members. The use of pivots connected to the mounting plate for latch or pawl members is illustrated by U.S. Pat. Nos. to H. A. Moline 2,970,325; B Roudebush 3,127,623; and W. W. Bennett 3,102,279. Also, the patent to R. I. Rumig, 2,866,210 requires a separate spring attached to the mounting plate for the latch member. These prior locking mechanisms for ski bindings are unnecessarily expensive and are, under some circumstances, difficult to operate during attachment and removal of the ski from the foot."}
{"text": "The present invention is directed to an improved flyback converter in which energy is re-circulated through the converter's flyback transformer windings in a controlled manner to eliminate the need for pre-load or bleeder resistors on the output(s) and to provide enhanced cross regulation on lightly loaded multiple output windings.\nConventional dc-to-dc power flyback converters designed for discontinuous mode operation subject the input power bus to a demand for current that ramps up from zero to a controlled peak level into the primary of a coupled inductor and is then interrupted. The input current then remains zero for an interval typically greater than one half of the total power period. At the instant the input current is interrupted, the stored energy in the coupled inductor forces current to immediately start flowing in the secondary or output winding. The wave shapes of the input current to the primary and the output current from the output winding are the same as that of a right triangle in each case. The current wave shapes thus produced are therefore the least desirable from the standpoint of stress management. In addition, during that period that follows the reset of the core as the current in the secondary falls to zero, and before the start of the next power cycle, the coupled inductor is performing no function.\nConventionally designed dc-to-dc power converters are typically provided with bleeder resistors that provide a minimum guaranteed load to prevent the output voltage from increasing at light load and to reduce the likelihood of a temporary spike on the output in the event of normal load interruption. The control loop bandwidth is limited and if a converter were operating at 50-150 kHz, a number of power cycles might be produced after the interruption of the load before the duty cycle on-time could be reduced, resulting in a jump in voltage on an unloaded output. Although necessary, a bleeder contributes to loss of efficiency. Flyback dc-to-dc power converters having multiple output windings typically have bleeder resistors on all outputs for the same reason as above and also because energy stored in the leakage inductance of the primary is coupled to the secondary windings and voltage on an isolated unloaded secondary output will continue to rise with the interruption of a normal load if there is no bleeder. Bleeders are also used on multiple outputs to enhance cross regulation."}
{"text": "Chicken eggs are one of the most important foods in the human diet, and are an exceptional source of proteins and fats, as well as amino acids and fatty acids. Every year in the United States an estimated 90 billion eggs are produced, with three fourths of these eggs being used for human consumption. An estimated 250 eggs per person are consumed annually in the United States.\nMany of the eggs consumed by humans are eaten as food ingredients, rather than directly as cooked eggs (such as boiled, fried, poached, etc.). In some cases whole eggs are used as food ingredients, for example as baking applications. However, it is often desirable to use just a portion of an egg as a food ingredient. For example, egg yolk is an excellent emulsifier and surfactant, and is an essential component of mayonnaise and various other foods. The egg yolk makes up approximately one third of the liquid weight of an egg, and is high in fats and fatty acids. Important fat soluble vitamins (A, D, E, and K) are found in egg yolk, as are unsaturated fatty acids (e.g. oleic acid, linoleic acid, palmitoleic acid, and linolenic acid) and saturated fatty acids (e.g. palmitic acid, stearic acid, and myristic acid). Egg yolks also contain some proteins, typically on the order of 2 to 3 grams out of about 15 to 20 grams of yolk within an egg weighing approximately 50 grams.\nThe egg white, known as well as the albumen, also has unique uses as a result of having high protein content. Egg whites are used in many products, such as to make mousse and to enhance protein content of foods. Egg white is approximately two-thirds of the total weight of an egg, with approximately 90 percent of that weight coming from water. The remaining weight of the egg white comes primarily from protein, along with various trace minerals, vitamins, some fats, and glucose. A typical large egg may contain 35 to 40 grams of egg white, of which about 4 to 5 grams are proteins. The most common protein in egg whites is ovalbumen, which accounts for over half of the proteins. Ovotransferrin and ovomucoid are additional primary proteins, with other proteins including ovoglobulin G2, Ovoglobulin G3, ovomucin, lysozyme, ovoinhibitor, ovoglycoprotein, flavoprotein, ovomacroglbulin, avidin, and crystatin. The egg white contains no dietary cholesterol, but does contain small quantities of other lipids and fats.\nThus, egg yolks are very high in fats, but low in proteins; while egg whites are very high in proteins, but low in fats. However, egg yolks do contain some proteins, and egg whites do contain small quantities of fats.\nDue to their different compositions and uses, it is often desirable to separate egg yolks and egg whites from one another. Various systems and methods have been developed for separation of eggs into yolks and whites. Separation of yolk and white from whole eggs can be done at high speeds under automated conditions, and can very effectively separate the yolks and eggs, with relatively little mixing of the yolk and eggs.\nDespite the uses of existing technology to separate yolks from whites, a need exists for further separation of egg components, including components in both the egg yolk and the egg white. This is true because the mere physical separation of the egg yolk from the egg white is not always sufficient to maximize use of the yolks and whites. For example, with regard to the yolk, it can desirable to also remove the yolk proteins from the yolk fats for at least two reasons: First, traditional uses of the yolk as an emulsifier improve upon reduction of the protein content. Second, the removal of the protein provides an isolated protein material that has further uses for applications where high-protein materials are desired and where specific yolk proteins are desired in isolated form. Similarly, with regard to the egg white, removal of non-protein materials creates a higher quality isolated protein. Also, separation of the proteins into different sizes and types can have significant benefits for production of specialized products.\nIn addition to the benefits associated with separating proteins and fats from nearly pure egg yolks and pure egg whites, a need also exists for separation of proteins and fats from mixtures that contain both egg yolks and egg whites. Such mixtures are created, for example, as a result of incidental breakage of yolks during the separation of the yolk from the white during the cracking process. Similarly, some egg white can remain with the egg yolk during cracking. These mixtures result in increased levels of proteins in the separated egg yolk, and increased levels of fats in the separated egg whites. The ability to separate the primary constituents (proteins and fats) within the mixtures can have meaningful advantages in terms of nutritional value and performance for specific applications (such as to create mayonnaise).\nYet another scenario for separation of fats and proteins in eggs arises due to production of eggs at hatcheries (or other facilities) where the eggs are not primarily raised for human consumption. For example, sterile eggs and non-incubated eggs from hatcheries are not produced or used for human consumption. Although these hatchery-derived eggs are typically not provided for regular human consumption, they are still of value as a source of fats and proteins. Currently the hatchery-derived eggs are not processed so as to be separated into yolk and white components, because they are produced in facilities that are not served by a high speed cracker and separator. Instead, the hatchery-derived eggs are often broken and the contents run through a separator to remove the egg shells, a process that mixes the egg yolks and proteins. This mixed egg yolk and white, which is considered inedible for human consumption, is typically used as a combined additive for uses such as animal feed. However, a need exists for a means for effectively separating the fats and proteins from these mixtures of inedible eggs so as to gain the maximum benefit of the proteins and fats so that they may be used in the most desirable manner. This is best achieved by separating them from one another.\nAs noted above, the most prevalent means of separating egg components is limited to the separation of whole yolks and whites from one another. This separation is commonly performed at high speeds and efficiencies using automatic equipment. However, alternative efforts for further separation of egg components into more isolated components have been attempted. Unfortunately, these methods have proven problematic for various reasons, including because the processes are inefficient, impractical, or have other deficiencies.\nAn example of such efforts is found in U.S. patent application Ser. No. 11/971,802 (“the '802 application”), assigned to Biova, Inc., and which is directed to a method of separating lipids from an egg mixture by using a cross-linking reagent. The cross-linking reagent is added to an egg mixture containing lipids and solubilized proteins, causing the lipids to crosslink so they can be separated from the proteins. Suitable crosslinking reagents include cyclobetadextran, silicon dioxide, colloidal silica material, fumed silica materials, and synthetic calcium silicate hydrates. The method of the '802 application may can include adjusting the pH level of the egg mixture to a pH at which the cross-linking reagent is functional so that cross-linking of the lipids occurs. The proteins are subsequently separated from the cross-linked lipids to provide a separated protein. The separated proteins may be obtained by subjecting the egg mixture to one or membranes or filters of various sizes to separate or further isolate proteins or populations of proteins of interest.\nAlthough it may be possible to separate proteins and fats derived from eggs using the teachings of the '802 application, significant drawbacks are associated with the methods and materials it discloses. First, the inclusion of the cross-linking reagent is problematic if the egg components will eventually be consumed by humans, in that there is often opposition on the part of the public and regulatory agencies to consume materials that have been subject to adulteration by organic reagents, such as cyclobetadextran, or by inorganic materials such as colloidal or fumed silica materials. The aversion to the use of these cross-linking reagents exists, in part, because of the potential for undesirable modification or contamination of the crosslinked ingredients. Although naturally occurring silica is widespread in nature, its use in food processing is unusual, and risks opposition by consumers even if little or none of the crosslinking agent remains in the separated components.\nSecond, the use of the crosslinking reagent necessarily adds expense to the processing of the egg mixture, both because of the cost of the crosslinking reagent itself, as well as the costs associated with the additional steps of crosslinking the fats and subsequent removal of the crosslinking agent (such as silica) from the fats after separation of the fats and proteins, assuming the fats are intended for further use. Thus, not only does the use of a reagent result in potentially problematic alteration and/or contamination of the egg material with crosslinker, especially the egg yolk (since it is the fats that are crosslinked, and fats are primarily found in the yolk), but the use of the reagent adds expense and complexity to the egg processing methods: A step must be made to add the reagent to the egg mixture, one to induce crosslinking of the fats, as well as steps to reverse the crosslinking, when possible, after separation of the crosslinked egg fats and egg proteins. These steps all take time, equipment, and effort.\nYet another problem with the methods taught in the '802 application is that the use of the reagents creates an issue with regard to waste silica (or other crosslinking agent), which must be removed from the fats for most uses of the fats, and which can result in creation of an undesirable waste stream that must be disposed of, even if it is not hazardous. Today there is an increased emphasis on processes that use limited resources and which produce little or no waste product, and the '802 application does not fully satisfy this objective.\nTherefore, a need exists for methods and equipment for separating egg components into proteins and fats (or other substituents, such as amino acids). Such methods and equipment should include the ability to separate the components of exclusively egg yolk materials, exclusively egg white materials, and mixtures of various levels of egg yolk and egg white. Preferably such methods and equipment can be efficient, cost effective, produced without undesirable alteration of the egg components (such as alteration with cross linkers), and do not create excessive undesirable waste streams."}
{"text": "1. Field of the Invention\nApparatuses and methods consistent with the present invention relate to a filtering circuit, and more particularly to a filtering circuit which can suppress a jammer in a wireless communication system.\n2. Description of the Related Art\nIn a multi-channel wireless communication environment, we have not only the signal in a desired channel, but also jammers in other channels. All of the jammers should be suppressed so strongly that the signal to noise ration (SNR) of the desired signal can be improved to a level necessary for successful wireless communication.\nA channel selection filter with the frequency characteristic to pass only the signal in the desired channel is needed to suppress jammers. The bandwidth of the filter corresponds to that of the channel bandwidth.\nA short-range communication system has a narrow channel bandwidth of less than 1 MHz, while some other communication systems have much wider channel bandwidths, for example, 20 MHz for Wi-Fi and more than 4 GHz for ultra-wideband (UWB).\nGenerally, in an integrated circuit (IC), the die area of an analog filter is inversely proportional to the frequency bandwidth. This means that the analog channel selection filter in a short-range communication system would typically occupy more than half of the entire wireless IC. A digital filter is used for channel selection in most of the commercial wireless IC's for a short-range communication system since a digital filter can be implemented with much smaller area than an analog filter.\nIn the case of using a digital filter for channel selection, analog signals including the desired signal and jammer signals must be converted to digital signals before suppressing jammers. Given that a wireless IC should provide successful communication under jammers 40 dB larger than the desired signal, an analog-to-digital converter (ADC) should have a 40 dB wider dynamic range than it should have without the presence of jammers. Therefore, the ADC is required to have higher resolution, leading to larger power consumption.\nThe block diagram shown in FIG. 1 is a proposed architecture to suppress jammers without any analog filters as an example of background art.\nThis architecture has two signal paths between a mixer circuit and an ADC.\nIn the first path, the jammers are extracted from the input signal by suppressing a desired signal. The input signal to the first path is converted to a digital signal using an ADC with a low resolution. Subsequently, only the desired signal is suppressed using a digital band stop filter. Finally, the digital signal that contains only the jammers is converted back to an analog signal.\nIn the second path, the input signal is delayed such that that jammers in the output of the second path are synchronous with jammers in the output signal of the first path.\nThe output signal from the first path consists of only jammers while the output signal from the second path consists of the desired signal and jammers. The jammers can be suppressed by subtracting the output signal in the first path from the output signal from the second path.\nSince this architecture shown in FIG. 1 does not contain analog filters, it can be realized with a smaller integrated circuit die area than previous architectures having analog filters. In addition, the required ADC resolution can be reduced since the jammers are suppressed before reaching the second ADC.\nHowever, the architecture shown in FIG. 1 has the severe problem of introducing noise to the system.\nThe delay circuit in the second path consists of sample and hold circuits connected in series as shown in FIG. 2.\nThe number of sample and hold circuits, Nc, is equal to the total delay time required for the delay circuit, Td, divided by the sampling time of the sample and hold circuit, Ts, as in the following equation 1:Nc=Td/Ts  Equation 1\nTd is set to a value equal to the delay time seen by the jammers going through the first path. This delay is nearly equal to the reciprocal of the bandstopwidth of the digital band stop filter in the first path, which corresponds with the bandwidth of the desired channel.\nTs is typically set to one-quarter of the reciprocal of the entire bandwidth of all the signals including the desired signal and jammers.\nThis architecture usually needs more than 100 sample and hold circuits in a narrow band communication system.\nSince a delay circuit introduces noise, such as thermal noise and switching noise, this architecture makes the SNR significantly worse."}
{"text": "The gain error, e.g., the difference between the actual transfer function and an ideal transfer function of a current feed-back instrumentation amplifier (CFIA) depends on the accuracy of matching between the CFIA's input transconductor and feedback transconductor. Some existing CFIAs try to address the gain error using an architecture in which each transconductor includes pre-amplifiers to create a transfer function proportional to the ratio between the degeneration resistors of the input and feedback transconductors, see FIG. 1. The advantage of this circuit architecture is that it provides a gain error that depends almost exclusively on matching of resistors which is superior to the matching of transistors. While the preamplifiers provide an overall low gain error, their noise adds to the total noise of the instrumentation amplifier, making this architecture less suitable for low noise designs. Another disadvantage is the presence of feedback loops around each pre-amplifier, which makes frequency compensation rather complex and difficult.\nIn another approach, the CFIA's transconductors are each based on a differential pair of transistors with degeneration. The noise is significantly lower than the CFIA implementation that uses preamplifiers, and there are no extra feedback loops. However, the gain error will also change: it will now be affected not only by the mismatch between the degeneration resistors of the input and feedback degeneration transconductors, but also by any mismatch between the degenerated differential paired transistors of each transconductor and by any mismatch between the respective tail currents being fed into each transconductor. As noted above, resistors match much better than transistors do, so the transistor mismatch will have a dominant effect on the gain error."}
{"text": "FIG. 1 diagrammatically illustrates a known stereoscopic display system comprising a stereoscopic display device in association with a complementary viewing device in the form of spectacles comprising two orthogonal polarisers intended to be disposed so that each polariser is positioned immediately in front of a respective eye of the viewer. The stereoscopic display device comprises a matrix of pixels half of which emit light polarised in a first direction and the other half of which emit light polarised in a second direction orthogonal to the first direction. For example the pixels may be arranged in columns and rows in the matrix, and alternate pixels or alternate rows or columns of pixels may be of the two different polarisation types. Thus the two eyes of the user see different images determined by switching of the pixels of the two polarisation types, and in this manner the viewer perceives a three dimensional scene. Ideally light pertaining to only one of the different images should be incident on each eye of the user over as broad a range of wavelengths as possible if the three-dimensional image is to be of the best quality.\nReferring to FIG. 1, such a known stereoscopic display device may comprise a layer of liquid crystal material 1 confined between two transparent electrodes 2 and 3 one of which is shown divided into two regions A and B denoting two pixels of the device. Furthermore the electrodes 2 and 3 and liquid crystal material 1 confined therebetween are disposed between two linear polarisers 4 and 5 having polarising axes whose relative orientation is determined by the mode of switching of the liquid crystal material in the device. For example, the polarising axes of the polarisers are normally aligned perpendicular to one another in a twisted nematic liquid crystal display (TNLCD) device. As is well known, the liquid crystal material of the device acts as a pixelated optical shutter enabling the emission of light to be controlled on a pixel-by-pixel basis, and accordingly it is not necessary to describe the switching operation of such a device in detail.\nIn the particular arrangement which is diagrammatically shown in FIG. 1 the polariser 5, which is positioned in front of the liquid crystal material 1 and associated electrodes 2, 3, has its polarisation axis 11 disposed vertically. Furthermore an optical retarder in the form of a quarter wave plate 6 is positioned in front of the polariser 5 and comprises two differently patterned regions A' and B' having optic axes 12 and 13 oriented respectively at -45.degree. and +45.degree. to the polarisation axis of the polariser 5. Thus light transmitted by the pixel A emerges from the region A' of the quarter wave plate 6 circularly polarised in a right-handed direction, as shown by the arrow 7, and light transmitted by the pixel B emerges from the region B' of the quarter wave plate 6 circularly polarised in a left-handed direction, as indicated by the arrow 8.\nIn order to view a stereoscopic image, it is necessary to ensure that the viewer's left eye is only aware of light transmitted by pixels of one type, such as the pixel A, whilst being permanently unaware of the state of pixels of the other type, such as the pixel B, whereas the viewer's right eye is only aware of light transmitted by pixels of the other type, such as the pixel B, whilst being permanently unaware of the state of pixels of the one type, such as the pixel A. To this end, viewing spectacles with different polarising arrangements are placed immediately in front of the viewer's eyes, only the arrangement in front of the viewer's right eye 14 being shown in FIG. 1 for the purpose of illustration.\nIn this case the arrangement comprises an optical retarder in the form of a quarter wave plate 9 having its optical axis 15 parallel to the optic axis 13 of the region B' of the quarter wave plate 6, and a polariser 10 having its polarisation axis 16 orientated at 90.degree. to the polarisation axis 11 of the polariser 5. Thus the region B' of the quarter wave plate 6 and the quarter wave plate 9 having parallel optic axis 13 and 15 behave as a half wave plate with its optic axis orientated at +45.degree. to the polarisation axis 11 and disposed between the crossed polarisers 5 and 10, thus allowing the transmission of light from the pixel B to the viewer's right eye 14. On the other hand, since the region A' of the quarter wave plate 6 and the quarter wave plate 9 have crossed optic axis and are disposed between the crossed polarisers 5 and 10, light transmitted from the pixel A is cancelled almost completely at all wavelengths (assuming the quarter wave plates 6 and 9 are made of the same material so that they have the same dispersion characteristics) so that substantially no light reaches the viewer's right eye 14 from the pixel A. Such fully achromatic behaviour due to summing of the effect of optical retarders to zero retardation is known, for example, from U.S. Pat. No. 4,719,507.\nAlthough not specifically shown in FIG. 1, it will be appreciated that the converse arrangement associated with the user's left eye comprises a quarter wave plate having its optic axis parallel to the optic axis 12 of the region A' of the quarter wave plate 6 and a polariser having its polarisation axis at 90.degree. to the polarisation axis 11 of the polariser 5. In this case such a quarter wave plate and the region A' of the quarter wave plate 6 having parallel optic axes behave as a half wave plate with its optic axis at -45.degree. between the crossed polarisers, and thus light transmitted from the pixel A reaches the viewer's left eye, whereas the crossed optic axes of the quarter wave plate and the region B' of the quarter wave plate 6 ensures that substantially no light from the pixel B reaches the viewer's left eye.\nHowever the two sequential quarter wave plates which form a half wave retarder (unlike a zero wave retarder) only act in the required manner at a single wavelength, so that their transmission properties between crossed polarisers are chromatic, rather than achromatic. This can be seen from the graph of FIG. 2 showing the transmission level as a function of wavelength for the light received by the viewer's right eye 14 from the region B', after passing through the analysing polariser 10, and from the region A' of the quarter wave plate 6. Whilst the light from the region A' is cancelled substantially completely at all wavelengths, the light from the region B', after passing through the analysing polariser 10, is substantially unattenuated at only one wavelength (about 525 nm) whereas the light at over wavelengths is attenuated to a lesser or greater extent. Of course, the light received by the viewer's left eye will have a similar transmission profile but with A' substituted for B' and B' substituted for A' in FIG. 2.\nInstead of having an achromatic dark state and a chromatic bright state as in the case of the graph of FIG. 2, it is possible for the device to be adapted to have a chromatic dark state and an achromatic bright state by providing a polariser 10 with a polarisation axis parallel to the polarisation axis 11 of the polariser 5. FIG. 3 is a graph of the transmission profile of the light received from the regions A' and B' of the quarter wave plate 6 by the viewer's right eye showing, as expected, the wavelength dependence of the light received from the region B'. Again the transmission profile of the light received by the viewer's left eye is similar with A' substituted for B' and B' substituted for A' in FIG. 3.\nFIG. 4 shows a graph illustrating the light levels received from the regions A' and B' of the quarter wave plate 6 by the viewer's right eye, but with polycarbonate retarders being used to form a quarter wave plate 9 in each part of the viewing spectacles. Such polycarbonate retarders are more conventional, and thus less costly, for such use. The performance of the system is not greatly compromised by the use of polycarbonate retarders.\nS. Pancharatnam \"Achromatic Combinations of Birefringent Plates\" Proceedings of Indian Academy of Sciences Vol. XLI, no. 4, Sec. A, 1955, pp. 130-136 and pp. 137-144 discuss the use of stacked uniform retarder films to improve achromaticity. However these papers are not specifically concerned with stereoscopic imaging systems."}
{"text": "For the design of digital circuits, designers often employ computer aided techniques. Standard languages, such as Hardware Description Languages (HDLs), have been developed to describe digital circuits to aid in the design and simulation of complex digital circuits. As device technology continues to advance, various product design tools have been developed to adapt HDLs for use with newer devices and design styles.\nAfter the HDL code is written and compiled, the design of an integrated circuit (IC) or a system which includes multiple ICs must be verified to be correct. Continually advancing processing technology and the corresponding explosion in design size and complexity have led to verification problems for complex circuit designs, such as Application Specific Integrated Circuits (ASICs) that are difficult to solve using traditional simulation tools and techniques.\nAs a result, some designers build prototype boards using multiple ICs such as field programmable gate arrays (FPGAs) to verify their ASIC designs. However, there are still problems with debugging the hardware design. When an error is detected during debug, designers may attempt to tap signals of interest from the circuit and use a logic analyzer to determine the cause of the error. However, this is a difficult process and is often not effective, especially in the case of intermittent errors. Errors that have already occurred are often difficult to repeat and reconstruct."}
{"text": "FIG. 1A illustrates a conventional, prior art H-bridge inverter circuit 110 for transferring power from a DC supply 120 to an AC load 115. The H-bridge inverter circuit 110 includes four switches 111-114. During the operation of the H-bridge inverter 110, the switches 111-114 turn on and off in pairs to deliver an AC current to the AC load 115. For example, at one instance, switches 111 and 114 are turned on, while the switches 112 and 113 are turned off. Current flows from DC power supply 120 through switch 111, to the AC load 115, through the switch 114, and to a ground voltage or a negative supply voltage. At another instance, switches 111 and 114 are turned off, switches 112 and 113 are turned on. During this instance, the current flows from the DC supply 120 through the switch 113, to the AC load 115, through the switch 112, and to the ground voltage or the negative supply voltage. By altering the direction of the current flow, the H-bridge inverter is configured to deliver an AC current from the DC supply 120 to the AC load 115.\nFIG. 1B illustrates a conventional, prior art sync bridge circuit 130 for transferring power from an AC supply 135 to a DC load 140. The sync bridge circuit 130 includes four switches 131-134. The sync bridge circuit 130 receives an AC current from the AC supply 135, and the switches 131-134 of the sync bridge circuit 130 turn on and off in pairs to deliver a DC current to the DC load 140."}
{"text": "This invention relates to hydraulic actuators for use in, e.g., robotic applications, and more particularly relates to force control of hydraulic actuators.\nAn actuator is generally defined as a device or mechanism that converts some form of energy into mechanical force or torque and linear or rotary velocity. A hydraulic actuator typically is connected to a high pressure fluid source and a flow control valve, e.g., a spool valve. Application of a small signal to the valve deflects the valve, allowing the fluid to flow, e.g., into one or more chambers driving a mechanical mechanism such as a piston provided in one or more of the chambers. With this action, the hydraulic actuator converts fluid flow into mechanical piston velocity, and provides the ability to control this velocity and corresponding mechanical position.\nHydraulic actuators are particularly well-suited for velocity and position control of robots and heavy equipment. Hydraulic systems also are generally characterized by the highest power density of modern controllable actuation systems because they are often operated at a pressure of as much as 3000 psi or greater. Hydraulic systems can also support large loads indefinitely while consuming minimal power. Given these attributes, hydraulic actuation systems are frequently the optimum choice for high force, high power density motion control applications such as automobile steering systems, airplane control surfaces, and heavy equipment operations employing, e.g., construction machinery.\nWhile hydraulic systems are in many respects optimal for velocity and position control, a number of inherent hydraulic system limitations constrain their applicability for force control. For most applications, force control requires an ability to sense and correspondingly control the forces of interaction between an actuator and the actuation environment. But in hydraulic systems, a measurement of the primary system variable, hydraulic pressure, does not fully enable such. Specifically, the pressure in a hydraulic chamber, e.g., a piston chamber, is not in general a good representation of the force at the actuator output. Hydraulic systems are in general very sensitive to contamination, such as foreign particles, in the hydraulic fluid. In order to limit such contamination, it is preferable to employ tight fluidic seals at the hydraulic piston and cylinder. Tight seals are found, however, to typically produce substantial stiction and coulomb friction during sliding, and to require a very high breakaway force, all of which contribute to force noise at the hydraulic actuator output and thereby limit the ability to accurately estimate output force. Dynamically, a range of factors, including non-linear flow characteristics, can be very difficult to control.\nThere have been attempts to reduce the sliding friction and stiction characteristic of tight hydraulic seals by, e.g., reducing the piston seal tolerance. In one example alternative, two or more sets of loose seals are employed, the first seal allowing leakage from the supply fluid pressure chamber and the second and following seals scavenging the leakage. Although this configuration can improve sliding characteristics, it is not cost effective for most applications and in practice can be very prone to leaks. As a result, for most applications only tight hydraulic seals can be employed.\nGiven this fundamental difficulty in estimating the output force of a hydraulic actuator as function of hydraulic pressure, hydraulic actuators have been largely limited to velocity and position control applications. Implementation of force control for robotic and other applications in a manner that exploits the high power density of hydraulic actuation has heretofore not been fully practical.\nThe invention provides the ability to effectively and precisely implement closed-loop force control of a hydraulic actuator, provided in accordance with the invention as a hydro-elastic actuator. The hydro-elastic actuator of the invention includes a hydraulic actuator, having a connection to hydraulic fluid and including a mechanical displacement member positioned to be mechanically displaced by fluid flow at the actuator. A valve is connected at the hydraulic actuator connection and has a port for input and output of fluid to and from the valve. At least one elastic element is provided in series with the mechanical displacement member of the hydraulic actuator and is positioned to deliver, to a load, force generated by the hydraulic actuator. A transducer is positioned to measure a physical parameter indicative of the force delivered by the elastic element and to generate a corresponding transducer signal. A force controller is connected between the transducer and the valve to control the valve, based on the transducer signal, for correspondingly actuating the hydraulic actuator and deflecting the elastic element.\nThe hydro-elastic actuator of the invention can be configured such that the force controller is connected to accept an input indicative of a desired actuator output force to be delivered to the load. Here the force controller is connected between the transducer and the valve to control the valve based on the transducer signal and the input, for correspondingly actuating the hydraulic actuator by an amount that delivers to the load the a desired actuator output force.\nThe hydro-elastic actuator of the invention provides the ability to make a high-fidelity measurement of the output force of a hydraulic system without measuring pressure or flow characteristics of the hydraulic system. The feedback control loop enables precise hydraulic system force control and control stability to a level not previously achievable without complicated control schemes to accommodate hydraulic characteristics. The high power and high force generation capabilities of the hydraulic actuator are preserved while providing shock tolerance and low system output impedance.\nThe hydro-elastic actuator of the invention is well-suited for an extremely broad range of applications, and is particularly effective at addressing high-force, high-power density applications. Robotics applications and heavy equipment operations, such as robotic fire fighting and earth moving, as well as telerobotic and haptic systems, are particularly well-addressed. Further, the important and growing class of biomimetic robots, and particularly dynamically-stable legged robots, which primarily rely on force control-based locomotion algorithms, are enabled by the invention to take on mass and scale not previously attainable.\nIn accordance with the invention, the hydraulic actuator can be provided as a hydraulic actuation chamber in which the mechanical displacement member is disposed with respect to the fluid connection. Here the fluid connection preferably consists of a fluid inlet and a fluid outlet of the chamber. This enables control of displacement of the displacement member by fluid flow into and out of the chamber. The valve can be connected to the fluid inlet and fluid outlet, and preferably is provided as a flow control valve. Whatever connection is employed between the valve and the fluid inlet and outlet, it preferably is dimensionally fixed. The valve port can include a connection for receiving fluid pumped by a fluidic pump.\nIn embodiments of the invention, the valve control signal is based on proportional or proportional-integral control of actuator output force. The valve control signal is in one embodiment an electrical current. This electrical control current is directed to the valve and is indicative of a controlled fluid flow to be produced through the valve. The electrical control current can be indicative of a controlled bi-state valve operation between a state of zero fluid flow and a state of maximum fluid flow through the valve. The force controller can further be connected to produce a fluid source control signal directed to a fluid source connected to the valve port. Here the fluid source control signal can be indicative of a controlled pulsed delivery of fluid to the valve in synchrony with the bi-state valve operation.\nIn accordance with the invention, the hydraulic actuator, when provided as a chamber, can consist of a linear piston chamber having a linear piston and a piston push rod extending out of the chamber, a rotary piston chamber having a rotary vane and a rotary shaft extending out of the chamber, or other suitable chamber configuration. When employed, a piston can be double- or single-acting. Preferably either a substantially non-leaky seal is provided at a location where the piston push rod or shaft extends out of the chamber or alternatively, at least one leaky seal and at least one leakage scavenger seal can here be employed.\nIn embodiments of the invention, the elastic element can be provided as a linear or a nonlinear elastic element. The elastic element can be provided specifically as at least one spring disposed in series with the hydraulic actuator displacement member, or as a plurality of springs positioned to together result in an elasticity provided in series with the hydraulic actuator displacement member. One or more coupling elements can be provided in series with and between the elastic element and the hydraulic actuator displacement member, as well as in series with and between the elastic element and the load.\nThe transducer can be provided as a potentiometer, a strain gauge, or other suitable sensing configuration, e.g., as a magnetic position sensor or an optical position sensor. In one embodiment, the transducer signal is based on deflection of the elastic element.\nThe hydro-elastic actuator of the invention can be produced of lightweight, low-cost, easily manufactured components. The elastic element force feedback control of the system preserves the high power and high force or torque generation of the system while providing precise force control and good force control stability. These characteristics are ideal for robotic and other mechanistic systems that interact with their environment. Other applications, features, and advantages of the invention will be apparent from the following description and associated drawings, and from the claims."}
{"text": "A typical door hinge for use on an entrance door has a pivot or knuckle from one side of which extends a leaf attached to the door jamb and from the other side of which extends a leaf fixed to the door so that the door can pivot relative to the jamb about a normally vertical axis. In DE 10 2004 016 769 one of the leaves is simply screwed to the door and the other leaf is formed as a comb and fitted to a holder in turn set in a housing recessed in the jamb. The holder can be shifted relative to the housing for adjustment of the hinge.\nIn line with known approaches, the housing is fastened to the door frame. The door hinge has a comb-like free hinge leaf that is inserted into the holder of the housing, which is fastened to the door jamb and is permanently fixed thereto. The holder has a fixing plate as well as a clamping plate and is guided on spindles. By turning the spindles, the position of the holder is modifiable relative to the housing that constitutes the mounting part. The constructive design allows the door to be adjusted from side to side, that is horizontally. In addition, the position of the hinge leaf can be corrected horizontally and vertically within the holder before the hinge leaf is secured in place in the desired position between the clamping plate and the fixing plate. The known housing requires a dimensionally stable support to mount the adjustment spindles. The fixing plate and clamping plate to secure the hinge leaf must be of solid construction to as to be able to hold the weight of the door. Construction of the housing is costly and requires an installation space with a depth of 22 to 25 mm.\nDE 199 21 796 describes a door hinge with a holder for receiving a hinge tab clamped between a tensioning plate and a clamping plate of the holder. For position adjustment of the holder, a spindle is provided that extends through a recess of the hinge leaf without rotational contact."}
{"text": "JP 3315993 B2 discloses a pant-type diaper comprising a liquid-pervious inner sheet, a liquid-impervious outer sheet and an absorbent structure sandwiched between these sheets. The diaper includes a front waist region adapted to lie on the wearer's ventral side and a rear waist region adapted to lie on the wearer's dorsal side. Transversely opposite side edges of the front waist region are joined to the associated side edges of the rear waist region so as to define a waist-opening and a pair of leg-openings. This diaper includes waist elastic members attached along a periphery of the waist-opening, and a waist-cuff extending, on the side of the diaper facing the wearer's skin, from the waist-opening toward the crotch region.\nIn this diaper, the waist-cuff is provided with an elastic member so that the waist-cuff may be put in close contact upon contraction of the elastic member so as to define a pocket between this waist-cuff and the diaper chassis and thereby to prevent body waste from leaking beyond the waist-opening.\nIn this diaper, the waist-cuff is left free from the diaper chassis except along the periphery of the waist-opening and at the transversely opposite ends of the front and rear waist regions. As a consequence, the pocket is apt to be deformed or even to be collapsed as the wearer's body moves, causing the body waste having been captured by the pocket to leak out from the diaper.\nJP 4215370 B2 discloses a further diaper in which the waist-cuff is fixed to the diaper chassis in the vicinity of transversely opposite side edges of an absorbent structure. In this diaper, therefore, the shape retention of the pocket is correspondingly improved so that the pocket is not easily collapsed.\nJP 4215370 B2 further discloses a diaper provided with leg cuffs formed along the transversely opposite side edges of the absorbent structure.\nWhen the absorbent structure absorbs or receives body waste, the absorbent structure is drawn downward under a weight of body waste absorbed or received by the absorbent structure. Therefore, the waist-cuff might be drawn apart from the wearer's skin under a weight of body waste absorbed or received by the absorbent structure when the waist-cuff is fixed to the chassis in the vicinity of the absorbent structure.\nWhen the diaper has a pair of leg-cuffs, the waist-cuff might partially overlap the leg-cuffs. In such diaper, if the waist-cuff is fixed in the vicinity of the side edges of the absorbent structure, the leg-cuffs may be collapsed to cause a leaking of body waste. Alternatively, the waist-cuffs might not overlap the leg-cuffs. In such diaper, there might be regions without both leg-cuffs and waist-cuff where leakage of body waste might occur."}
{"text": "1. Field of the Invention\nThe present invention relates to apparatuses and methods for forming workpieces, and particularly to apparatuses and methods that readily and accurately perform a variety of forming operations on metallic workpieces.\n2. Related Art\nWith incessant development of modern manufacturing industries, more and more components having a myriad of different shapes are being required. Shapes commonly found on such components include, for example, inclined surfaces, arcuate surfaces, and steps. A variety of apparatuses and methods are available for forming such components.\nConventional apparatuses and methods for forming workpieces only allow the workpiece to be machined linearly. Arcuate surfaces are often formed using electrodischarge machines. However, electrodischarge machine is complicated and time-consuming. Furthermore, a special electrode is required for every different workpiece. This is unduly costly. Moreover, it is difficult to control the precision with which an electrodischarge machine forms a workpiece. Thus quality control of the workpieces is problematic.\nAn improved apparatus and method for forming workpieces which overcomes the above-mentioned problems and shortcomings is desired.\nAccordingly, an object of the present invention is to provide an apparatus which can readily and accurately form a surface having three dimensions on a metallic workpiece.\nAnother object of the present invention is to provide a method which can efficiently and accurately form a surface having three dimensions on a metallic workpiece.\nTo achieve the above-mentioned objects, a forming apparatus in accordance with the present invention comprises a workbench and a pivot device. The workbench comprises a magnetic base for holding a workpiece thereon. The pivot device comprises a pivot plate and a fixed plate. The pivot plate is fixedly attached below the base, and rotatably seated on the fixed plate. The fixed plate defines a pair of cutouts in two corners thereof respectively, to form an arched-shaped top portion. A pair of locking plates is attached on a first side wall of the pivot plate, and another locking plate is attached on a second side wall of the pivot plate. When the locking plates are oriented vertically, bottom extremities of the locking plates are located between a top extremity of the top portion and a bottom extremity of the top portion. The locking plates can therefore pass over the two corners corresponding to the cutouts, but cannot pass over other corners of the fixed plate. When the locking plates are oriented horizontally, the bottom extremities of the locking plates are higher than the top extremity of the top portion. The locking plates can therefore pass over every corner of the fixed plate. By adjusting orientations of the locking plates, the combined base and pivot plate can be rotated within a variety of predetermined different angles.\nOther objects, advantages and novel features of the present invention will be drawn from the following detailed description of preferred embodiments of the present invention with attached drawings, in which:"}
{"text": "1. FIELD OF THE INVENTION\nThe present invention relates to a focus condition detecting the device for a camera which detects focus condition of an objective lens of a camera from an object light having passed through the object lens.\n2. DESCRIPTION OF THE PRIOR ART\nThere has been known a focus condition detecting device in which two images are formed by refocusing light bundles of an object light having passed through first and second areas of an objective lens which symmetric with respect to the optical axis of the objective lens. A relative distance between these two images is calculated and a defocus amount of the focus position from a predetermined focus position and a direction of the defocus are determined based on the relative distance calculated.\nA typical optical system for such a focus condition detecting device as mentioned above is shown in FIG. 1.\nIn this system, a condenser lens 6 is usually arranged on a predetermined focal plane 4 of an objective lens 2 or on a plane positioned behind the focal plane and two refocusing lenses 8 and 10 are arranged rearside of the condenser lens 6. There are also arranged two image sensors 12 and 14 on each of focal planes of two refocusing lenses. Each of the image sensors 12 and 14 is comprised of a CCD (charge coupled device) image sensor having a plurality of light sensing elements.\nThe condenser lens 6 is so designed that effective ranges of the first and second refocusing lenses 8 and 10 can be projected within the exit pupil of the objective lens in order to form first and second images by light bundles having passed through the exit pupil and first and second refocusing lenses.\nAs shown in FIG. 2 schematically, when an image (A) of an object is focused frontside of the predetermined focal plane 4, two images a and a' are formed on the image sensors 12 and 14 so as to approach to each other with respect to the optical axis 18 of the objective lens. On the contrary to the above, when an image (B) of an object is focused rearside of the predetermined focal plane 4, two images b, b' are formed apart from each other. If an image is focused just on the predetermined focal plane 4, a distance between two corresponding points of the two images formed on two image sensors 12 and 14 becomes a specific value which is determined by the composition of the optical system of the focus condition detecting device. Accordingly, a focus condition of the objective lens can be determined from the distance between two corresponding points of the two images formed on the image sensors 12 and 14.\nThe detection of the distance mentioned above can be made as follows.\nFIG. 3 shows a principle of operation therefor schematically. As shown therein, the first and second image sensors 12 and 14 are comprised of ten and sixteen cells of photodiodes (a.sub.1, . . . , a.sub.10) and (b.sub.1, . . . , b.sub.16), respectively. Now, assume that each symbol of an individual cell represents each output thereof. Considering now sets each of which is comprised of ten successive cells included in the second image sensor, seven sets B.sub.1, B.sub.2, . . . B.sub.7 can be obtained as shown in FIG. 3. The focus condition can be sought by calculating an individual correlation relation between the image received by ten cells of the first image sensor 12 and the image received by each of seven sets of the second image sensor 14.\nNamely, correlation calculations are made with use of correlation functions: ##EQU1##\nFor example, if the image detected by the first sensor 12(A.sub.1) coincides with the image detected by the first set B.sub.1 of the second image sensor 14, the correlation function S.sub.1 becomes minimum among seven correlation functions S.sub.1, S.sub.2, . . . , S.sub.7. When either one set of the second sensor is found which gives the minimum value with respect to these correlation functions, the distance between two images is determined from the number of the set having been found minimum and a focus condition is detected based on the distance determined. These calculations are carried out by a correlation calculation means 16.\nIn FIG. 4 showing one scene viewed through a view-finder of a camera, the focus detection area A is limited to a relatively narrow center area of the field of view of the objective lens. In this case, the image of one person (main object) existing near to the camera and the background image comprised of trees and a hill ar entered into the focus detection area A and, due to this situation, the focus condition detecting device is liable to detect an intermediate distance between the person and the background as a correct distance of the object.\nMeanwhile, in the case that a few objects having different distances from the camera respectively are entered into the focus detection area, the main object usually exists at a position nearer than those of the sub objects (background) and, usually, may be one or more persons. And the sub object is usually a background of the person. Considering such an ordinary or typical object as mentioned above, the image of the main object such as a person or persons has a contrast higher than that of the background usually.\nIn Japanese patent laid open publication No. 126517/1984, there is proposed a focus detecting device in which the focus detection area is divided into a plurality of blocks and a focus detection operation is carried out in every block and the credibility about individual focus detecting operation in every block is checked and the result obtained through the focus detecting operation having the highest credibility is employed as focus detection data.\nHowever, according to this focus detecting device, a long calculation time is needed sine the credibility of focus condition detection has to be checked at every block divided."}
{"text": "1. Field of the Invention\nThe present invention relates to a racing type pedal, and more particularly to a racing type pedal for a bicycle.\n2. Description of the Related Art\nIn general, the racing type pedal for a bicycle includes a road-use pedal and a mountaineering pedal, and the two different types of pedals need to mate with two different types of shoes and snap blocks. In practice, the snap block is mounted on the bottom of the bicycle shoe and snapped onto the bicycle pedal to secure the bicycle shoe on the bicycle pedal.\nA conventional racing type pedal 3 in accordance with the prior art shown in FIGS. 10 and 11 comprises a main body 30, a hollow pivot hook 31 pivotally mounted on the main body 30 by a pivot shaft 32 and having a side formed with a vertical plate 310 having a through hole 3101, an adjusting bolt 34 mounted on the plate 310 of the pivot hook 31 extended through the through hole 3101 of the plate 310 of the pivot hook 31, a press block 35 mounted on the distal end of the adjusting bolt 34, and two torsion springs 33 each mounted on the pivot shaft 32 and each having a first end 330 urged on a block 301 of the main body 30 and a second end 331 urged on the press block 35.\nHowever, when the adjusting bolt 34 is rotated to move the press block 35 to adjust the torque of the pivot hook 31, the adjusting bolt 34 and the press block 35 are subjected to a shearing force applied by the torsion springs 33, so that the adjusting bolt 34 is easily deflected from the press block 35 during movement, thereby disengaging or wearing the adjusting bolt 34 and the press block 35. In addition, the plate 310 of the pivot hook 31 is subjected to a force applied by the adjusting bolt 34, so that the plate 310 of the pivot hook 31 is easily distorted or deformed due to the applied force during a long-term utilization."}
{"text": "Portable, mobile computers such as slate computers have been used in business. The mobile computers are likely to be lost or stolen. Some of the mobile computers, therefore, have a function of deleting data in a terminal by remote control when it is lost or stolen (hereinafter, referred to as a remote wipe).\nA remote wipe performed when a mobile computer is lost or stolen has a problem that the probability of success is low for such a reason that the computer is not connected to a network, and whether or not the remote wipe has succeeded cannot be determined."}
{"text": "The invention relates to a line-voltage-operated incandescent lamp having\na translucent lamp vessel sealed in a vacuum-tight manner; PA0 current supply conductors entering through the wall of the lamp vessel and mechanically interconnected inside the lamp vessel by means of glass beams with which they constitute a frame; PA0 a filament comprising several helically wound sections which are stretched at least substantially parallel to each other in the frame, each end of the filament being electrically connected to a respective current supply conductor; and PA0 two sets of supports which are anchored in a respective glass beam of the frame. Each support consists of a wire terminating in a hook and each hook engages the filament between a respective pair of adjacent sections of the filament.\nSuch a lamp is known from British Patent Specification 2,069,233 to which U.S. Pat. No. 4,367,428 corresponds.\nThe known lamp has a frame of a rigid construction, as a result of which the lamp is capable of withstanding shocks. In this respect, the known lamp is superior to lamps which were previously commercially available, in which two tungsten wires are anchored each with one end in one of the glass beams of the frame, while the other end of these wires is helically wound and grips around a respective current supply conductor and is fixed thereto. In these commercially available lamps, the length of the frame and hence the stress in the filament can be readily adjusted by slipping the helically wound ends further or less far over the current conductors, but the frame is not very capable of withstanding shocks.\nIn the manufacture of the lamp according to the that British Patent Specification, it has proved to be difficult to mount the filament with the correct tension in all its sections."}
{"text": "In sports that utilize equipment having an elongated shaft, the shaft must ideally be both lightweight and strong. However, these two requirements are often incompatible, in that reduction in weight often may cause a loss of strength and vice versa. Ideally, a shaft should have sufficient strength to withstand the stresses and deformation that arise during use and the impacts that it may be subjected to during play. This is particularly true in contact sports such as ice hockey, field hockey, lacrosse, ringuette, and others. Ideally, the elongated shafts used in those sports, must be able to withstand a large number of impacts, which impacts may often be concentrated at the edges, i.e. the corners or angles thereof formed by the meeting of two adjacent lateral sides of the shaft which, over time, may result in increased damage to the structure of the shaft and ultimately, premature failure thereof. The same concerns apply to the blade of a stick, which is subjected to many impacts, particularly on the upper surface.\nHockey sticks (including goalie sticks), field hockey sticks, lacrosse sticks, ringuette sticks and other such sports sticks may have shafts which may be made from a variety of materials including wood, aluminum, plastic and composite materials such as fiberglass, graphite and Kevlar™ or a combination of any of them. Some shafts are full (i.e not hollow), while others comprise four (relatively) thin side walls forming a peripheral box having a hollow core. Most blades are full, i.e. not hollow. Alternatively, some shafts and some blades have a composite construction having various layers of materials sandwiched as a core. Materials are usually selected for their physical properties in an attempt to improve performance, longevity, etc. . . . Some composite shafts may have lower durability, but are still popular because of their light weight and superior stiffness. Wood shafts are cheap, but are not especially light, stiff or durable while aluminum shafts can have a relatively short life as they are prone to bending failure. Cost is often a criterion in material selection. All of these shafts may be particularly vulnerable to failure along their edges, i.e. where one side surface intersects with an adjacent side surface, often at 90°. Impacts are often concentrated at these edges, precisely where there is less material to absorb and dissipate said impacts. The same problem is experienced by the blades. Lastly, sticks that are the subject of repeated impact on their edges rapidly become worn and tired-looking, with paint and decals worn off, and nicks and gouges therein. Some players do not like their equipment looking shabby.\nThere is therefore a need for a sports apparatus shaft that has an increased ability to withstand impact along its edges.\nAccordingly, it is an object of the present invention to provide a sports shaft where there is provided added protection at the edges thereof.\nIt is another object of the present invention to provide a sports shaft where there is removed some material along at least one longitudinal edge thereof, which material is replaced with another material more suited for absorbing and resisting impacts.\nIt is another object of the present invention to provide a sports shaft where the shape of the shaft is such that grooves are provided along at least one of the edges so as to provide a volume to be filled with a material more suited for absorbing and resisting impact.\nIt is another object of the present invention to provide a blade shaft wherein at least one groove is provided along the upper edge so as to provide a volume to be filled with a material more suited for absorbing and resisting impact.\nAn additional object of the present invention is to provide a sports shaft wherein bumpers are selectively provided on the edges thereof to absorb and distribute the shock of an external impact, ideally in a direction perpendicular to the line of action of impact.\nAn additional object of the present invention is to provide a sports shaft wherein said bumpers are made from an elastomeric material."}
{"text": "1. Field of the Invention\nThe present invention relates generally to file systems, and in particular, to a method, apparatus, and article of manufacture for solving the difficulty incurred by a program on one computer (the client) being completely reliant on the internal file structure of another computer (the server) in order to access the server's files.\n2. Description of the Related Art\nComputer use has proliferated and continues to expand as technological capabilities rapidly increase. With this proliferation, the development and use of different computers running different types of operating system having different internal file systems has also increased. In a network environment, the different computers may attempt to communicate with each other including transmitting/receiving information and services to/from each other. However, when a computer program attempts to access a file on any other particular computer, the computer program is subjected to the quirks of the file system upon that computer. This includes the fact that the program trying to access the file needs to know both the name of the file, and where the file actually resides internally on the computer. Such information may not be known or may be difficult or impossible to learn and/or use. Accordingly, a method is needed for accessing information on files without knowledge of a computer's underlying file structure. These problems may be better understood by describing prior art file systems and network communications.\nDepending on the type of operating system a computer is executing, the internal file system varies considerably. For example, if a computer is running the Windows™ operating system, the file system may be composed of one or more drives; either local within the computer or logically mapped to a network connection. Each drive is assigned a letter, starting from C (A and B are reserved) and ending at Z, followed by a colon (:) character. Within a given drive exists a hierarchical directory structure, using variable length names and the backward slash (\\) character to indicate a level in the hierarchy. Examples of a hierarchical directory structure in the Windows operating system include those illustrated in Table 1:\nTABLE 1C:\\TempA directory name called Temp located onthe C drive.C:\\Temp\\FredA directory called Temp further containinga sub-directory called Fred.C:\\Temp\\Fred\\file.txtA file residing in the Fred sub-directory ofthe Temp directory.\nIf a computer is running the UNDX™ operating system, the file system is composed of one or more “mount” points; either local within the computer or logically mapped to a network connection. Each mount point is assigned a name, and is referenced using two forward slashes (//). Within a given mount point exists a hierarchical directory structure, using variable length names and the forward slash (/) character to indicate a level in the hierarchy. Examples of a hierarchical directory structure in the UNIX™ operating system include those in Table 2:\nTABLE 2//TempA directory name called Temp located onthe default mount point.//Temp/FredA directory called Temp further containinga sub-directory called Fred.//Temp/Fred/file.txtA file residing in the Fred sub-directory ofthe Temp directory.\nIf a computer is running the OS/390™ operating system, the file system is composed of one or more high-level prefix names that are controlled by a central catalog. Beneath the high level prefix lies one or more 1 to 8 character names, separated by a period (.) character, with a total length not exceeding 44 characters.\nThe combination of high-level prefix and the period-delimited suffix is called a “dataset name”. There are many types of datasets on OS/390. Each of the datasets further complicates the indicated format. One of the more common datasets is called a Partition Dataset or PDS. A PDS may contain one ore more files, indicated by use of open and close parentheses ( ). Examples of a hierarchical data structure in the OS/390 operating system include those illustrated in Table 3:\nTABLE 3HILEV.temp.fred.FileA dataset containing many files.HILEV.temp.fred.File(txt)A dataset name calledHILEV.temp.fred.file containing a filecalled txt.\nThe details of the file systems are not necessary for understanding the present invention. However, it should be recognized that the file systems are all different. In this regard, a computer program trying to access a file on any particular computer is subjected to the quirks and details of the file system upon that computer. Such details include the fact that the program trying to access the file needs to know both the name of the file, and where the file actually resides internally on the other computer.\nThe prior art necessity for knowledge of the file system details may be illustrated by the following example. Assume three different computers (servers) are communicating: one running OS/390, one running Windows™, and the third one running UNIX. Each server houses a copy of a file. In each case, the file is the same in the sense that its content is identical. It is only the file system upon which the file resides that varies from one machine to the next.\nA fourth computer, acting as a client, wishes to retrieve this file from any server. The client does not care which server it retrieves the file from, since the file will be the same no matter which server handles the request. Thus when the client connects with the server, it wishes to merely “ask for the file” without the need to care about the server's file system. To paraphrase in English, the client request may be “Excuse me server, could you please give me file X?” as distinct from “Excuse me server, could you please give me file X which you have located in quadrant Y of partition Z on drive Q?”\nWhen considering the second request above, it should be clear that the client would not know that this particular server uses the term “quadrant” while a different server may use the term “directory” while the third server may use the term “fandangle”. Furthermore, even if the client could be made to be completely sensitive to the server's file system and know the proper form for requesting a file (e.g., the client speaks in the vocabulary of the server), the client would not know where the file actually was located. In other words, the client request may be “Excuse me server number 2, I speak your language, so could you please refer to drive C, directory Alpha, delimiter ‘\\’, sub-directory Beta, and send me a copy of file X?” In response, the server may return the file to the client. However, to personify the server with human emotion, the server may respond: “Hello Client, yes here is file X, but how did you know where I had it located?”\nIn view of the above, it may be understood that a program trying to access a file in the prior art needs to know both the name of the file and where the file actually resides internally on the other computer. Gaining and maintaining such knowledge is problematic."}
{"text": "Conventionally, as a means of displaying images, projection-type optical display apparatuses such as projectors are known. Such optical display apparatuses require an optical illumination apparatus for efficiently and uniformly illuminating the optical image formed on a display panel, such as a reflective liquid crystal display panel. FIG. 10 is a diagram conceptually showing an example of the construction of an optical display apparatus employing a conventional optical illumination apparatus.\nIn FIG. 10, reference numeral 101 represents a light source, and reference numeral 102 represents a reflector disposed so as to partially surround the light source 101. A PBS (polarizing beam splitter) prism unit 103 is disposed immediately behind the reflector 102, i.e., on the right side thereof in FIG. 10. The PBS prism unit 103 includes a plurality of PBS prisms arranged parallel to one another. The PBS prism unit 103 splits the light from the light source 101 into two differently polarized types of light. Of the individual PBS prisms 103a and 103b, those which let out S-polarized light as described later have half-wave plates 104 disposed immediately behind them.\nBehind the PBS prism unit 103 (i.e., on the right side thereof in FIG. 10) are disposed, in order of arrangement, a first lens array 105, then somewhat away therefrom, a second lens array 106, and a superimposing lens 107 immediately behind it. The first lens array 105 has a plurality of lens cells 105a arranged in a rectangular, grid-like array having an aspect ratio substantially identical to that of a display panel 109 to be described later. Similarly, the second lens array 106 also has a plurality of lens cells 106a arranged in a rectangular, grid-like array. However, the shape of the lens cells 106a of the second lens array 106 is not necessarily geometrically similar to that of the lens cells 105a. \nThe images from the individual lens cells 105a of the first lens array 105 are, by the second lens array 106 and the superimposing lens 107 disposed immediately behind it, superimposed on one another in the vicinity of the focal point of the superimposing lens 107. The display panel 109 is disposed at the focal point of the superimposing lens 107. The display panel 109 is illuminated in a telecentric fashion by a condenser lens 108 disposed immediately in front of it. The components from the first lens array 105 through the superimposing lens 107 mentioned above together constitute an optical integrator system. It is to be noted that, in all the diagrams referred to in the present specification, irrespective of whether they relate to prior-art examples or to embodiments of the present invention, light beams are represented by their optical axes alone.\nThe light emitted from the light source 101 is reflected from the reflector 102, and is thereby formed into a substantially parallel beam and directed to the PBS prisms 103a of the PBS prism unit 103. Here, P-polarized light, indicated by solid lines P, is transmitted straight through the PBS prisms 103a. On the other hand, S-polarized light, indicated by broken lines S, is reflected inside the PBS prisms 103a so as to be directed to the outwardly contiguous PBS prisms 103b, and is then reflected again inside the PBS prisms 103b so as to exit therefrom, still as S-polarized light. That is, by the PBS prism unit 103, the light from the light source 101 is split into two differently polarized types of light in the direction of the longer sides of the display panel 109, i.e., in a vertical direction along the plane of the figure.\nThe S-polarized light exiting from the PBS prisms 103b is transmitted through the half-wave plates 104 disposed immediately behind the PBS prisms 103b and is thereby converted into P-polarized light. That is, a portion of the light from the light source 101 has its polarization converted first by the PBS prisms 103b of the PBS prism unit 103 and then by the half-wave plates 104, and eventually comes out as uniformly P-polarized light. This arrangement constitutes a polarization conversion device. Here, the type of light into which the light from the light source 101 is converted does not necessarily have to be P-polarized light, but can be of other polarizations. The arrangement described thus far, starting with the light source 101 and ending immediately in front of the display panel 109, constitutes an optical illumination apparatus.\nThe light thus converted into uniformly P-polarized light is then directed through the above-mentioned optical integrator system to the display panel 109. The display panel 109 modulates, pixel by pixel, the light it is illuminated with according to the display data fed thereto, and emits the modulated light. The light thus emitted then enters an optical projection system 110. The display data presented on the display panel 109 is projected, as an image, onto a screen (not shown) through this optical projection system 110. Reference numeral 110a represents an aperture stop disposed in the optical projection system 110.\nFIG. 11 is a diagram conceptually showing another example of the construction of an optical display apparatus employing a conventional optical illumination apparatus. In this figure, reference numeral 201 represents a light source, and reference numeral 202 represents a reflector disposed so as to partially surround the light source 201. Behind the reflector 202 (i.e., on the right side thereof in FIG. 11) are disposed, in order of arrangement, a first lens array 203 and, then somewhat away therefrom, a second lens array 204. The first lens array 203 has a plurality of lens cells 203a arranged in a rectangular, grid-like array having an aspect ratio substantially identical to that of a display panel 209 to be described later. Similarly, the second lens array 204 also has a plurality of lens cells 204a arranged in a rectangular, grid-like array. However, the shape of the lens cells 204a of the second lens array 204 is not necessarily geometrically similar to that of the lens cells 203a. \nA PBS (polarizing beam splitter) prism array 205 is disposed immediately behind the second lens array 204. The PBS prism array 205 includes a plurality of PBS prisms arranged in an array. The PBS prism array 205 splits the light from the light source 201 into two differently polarized types of light. Of the individual PBS prisms 205a and 205b, those which let out S-polarized light as described later have half-wave plates 206 disposed immediately behind them.\nA superimposing lens 207 is disposed behind the PBS prism array 205. The images of the individual lens cells 203a of the first lens array 203 are, by the second lens array 204 and the superimposing lens 207, superimposed on one another in the vicinity of the focal point of the superimposing lens 207. The display panel 209 is disposed at the focal point of the superimposing lens 207. The display panel 209 is illuminated in a telecentric fashion by a condenser lens 208 disposed immediately in front of it. The first lens array 203, the second lens array 204, and the superimposing lens 207 mentioned above together constitute an optical integrator system.\nThe light emitted from the light source 201 is reflected from the reflector 202, and is thereby formed into a substantially parallel beam and passed through the first lens array 203 and the second lens array 204, so that the light exiting from the individual lens cells 204a of the second lens array 204 enters corresponding ones of the PBS prisms 205a of the PBS prism array 205. Here, P-polarized light, indicated by solid lines P, is transmitted straight through the PBS prisms 205a. On the other hand, S-polarized light, indicated by broken lines S, is reflected inside the PBS prisms 205a so as to be directed to the contiguous PBS prisms 205b, and is then reflected again inside the PBS prisms 205b so as to exit therefrom, still as S-polarized light.\nThe S-polarized light exiting from the PBS prisms 205b is then transmitted through the half-wave plates 206 disposed immediately behind the PBS prisms 205b and is thereby converted into P-polarized light. That is, a portion of the light from the light source 201 has its polarization converted first by the PBS prisms 205b of the PBS prism array 205 and then by the half-wave plates 206, and eventually comes out as uniformly P-polarized light. This arrangement constitutes a polarization conversion device. Here, the type of light into which the light from the light source 201 is converted does not necessarily have to be P-polarized light, but can be of other polarizations. The arrangement described thus far, starting with the light source 201 and ending immediately in front of the display panel 209, constitutes an optical illumination apparatus.\nThe light thus converted into uniformly P-polarized light is then directed through the superimposing lens 207 to the display panel 209. The display panel 209 modulates, pixel by pixel, the light it is illuminated with according to the display data fed thereto, and emits the modulated light. The light thus emitted then enters an optical projection system 210. The display data presented on the display panel 209 is projected, as an image, onto a screen (not shown) through this optical projection system 210. Reference numeral 210a represents an aperture stop disposed in the optical projection system 210.\nIn the conventional optical illumination apparatus constructed as shown in FIG. 10, polarization conversion is performed immediately behind the light source 101. Therefore, the light emitted from the light source 101 and then reflected from the reflector 102 has its beam diameter enlarged to about twice its original beam diameter as a result of the polarization conversion. This diminishes the f-number of the illumination light Ia that strikes the display panel 109 and thus diminishes the f-number of the projection light Ea that emanates from the display panel 109, making the burden on the optical projection system 110 heavier.\nOn the other hand, in the conventional optical illumination apparatus constructed as shown in FIG. 11, the light emitted from the light source 201 and then reflected from the reflector 202 experiences no enlargement of its beam diameter. Therefore, no diminishing occurs in the f-number of the illumination light Ib that strikes the display panel 209 nor in the f-number of the projection light Eb that emanates from the display panel 209. Thus, no extra burden is placed on the optical projection system 210. However, the light from the light source 201 is not converted into uniformly polarized light until it has passed through the second lens array 204. Therefore, in this optical illumination apparatus, no space is available for inserting a polarization-dependent color switching device such as those used in the embodiments of the present invention to be described later. That is, this optical illumination apparatus does not permit a so-called color sequential illumination method using such a color switching device."}
{"text": "The present invention relates to a system and method for the treatment of disorders of the vasculature. More specifically, a system and method for treatment of abdominal aortic aneurysm and the like, which is a condition manifested by expansion and weakening of the aorta below the diaphragm. Such conditions require intervention due to the severity of the sequelae, which frequently is death. Prior methods of treating aortic aneurysm have consisted of invasive surgical methods with graft placement within the aorta as a reinforcing member of the artery. However, such a procedure requires a surgical cut down to access the vessel, which in turn can result in a catastrophic rupture of the aneurysm due to the decreased external pressure from the organs and tissues surrounding the aorta, which are moved during the procedure to gain access to the vessel. Accordingly, surgical procedures have a high mortality rate due to the possibility of the rupture discussed above in addition to other factors. Other factors can include poor physical condition of the patient due to blood loss, anuria, and low blood pressure associated with the aortic abdominal aneurysm. An example of a surgical procedure is described in a book entitled Surgical Treatment of Aortic Aneurysms by Denton A. Cooley, M.D., published in 1986 by W.B. Saunders Company.\nDue to the inherent risks and complexities of surgical procedures, various attempts have been made in the development of alternative methods for deployment of grafts within aortic aneurysms. One such method is the non-invasive technique of percutaneous delivery by a catheter-based system. Such a method is described in Lawrence, Jr. et al. in “Percutaneous endovascular graft: experimental evaluation”, Radiology (May 1987). Lawrence described therein the use of a Gianturco stent as disclosed in U.S. Pat. No. 4,580,568. The stent is used to position a Dacron fabric graft within the vessel. The Dacron graft is compressed within the catheter and then deployed within the vessel to be treated. A similar procedure has also been described by Mirich et al. in “Percutaneously placed endovascular grafts for aortic aneurysms: feasibility study,” Radiology (March 1989). Mirich describes therein a self-expanding metallic structure covered by a nylon fabric, with said structure being anchored by barbs at the proximal and distal ends.\nOne of the primary deficiencies of the existing percutaneous devices and methods has been that the grafts and the delivery catheters used to deliver the grafts are relatively large in profile, often up to 24 French and greater, and stiff in bending. The large profile and bending stiffness makes delivery through the irregular and tortuous arteries of diseased vessels difficult and risky. In particular, the iliac arteries are often too narrow or irregular for the passage of a percutaneous device. Because of this, non-invasive percutaneous graft delivery for treatment of aortic aneurysm is not available to many patients who would benefit from it.\nAnother contraindication for current percutaneous grafting methods and devices is a vessel treatment site with high neck angulation which precludes a proper fit between the graft and the vessel wall. An improper fit or seal between the graft and the vessel wall can result in leaks or areas of high stress imposed upon the diseased vessel which lead to reduced graft efficacy and possibly rupture of the aneurysm.\nWhile the above methods have shown some promise with regard to treating abdominal aortic aneurysms with non-invasive methods, there remains a need for an endovascular graft system which can be deployed percutaneously in a small diameter flexible catheter system. In addition, there is a need for a graft which conforms more closely to the contours of an aortic aneurysm which are often quite irregular and angulated and vary from patient to patient. The present invention satisfies these and other needs."}
{"text": "Roku, Inc., of Saratoga, Calif., the assignee of the present invention, currently offers a streaming media platform for connection to a television. In particular, the streaming media platform utilizes a cable (i.e., a High Definition Multimedia Interface or HDMI cable) to connect to external ports of a television. The streaming media platform uses a wireless connection to access the internet. Consequently, streamed media from the internet is wirelessly collected by the streaming media platform and is applied to the cable for display on a television.\nThere are commercially available televisions with resident internet access circuits. These internet access circuits may access the internet through a wireless connection or a wired connection, such as an Ethernet cable connection. Many purchasers are reluctant to purchase a television with a built-in internet access circuit because such purchasers do not want to be committed to a certain internet connection interface."}
{"text": "1. Field of the Invention\nThe invention relates to the field of optical fibers for wavelength division multiplex transmission networks.\n2. Description of the Prior Art\nIncreasing information bit rates on this type of network impose compensation of chromatic dispersion and dispersion slope over bands of the spectrum other than the C band. The S band extends approximately from 1 460 nm to 1 530 nm. The C band extends approximately from 1 530 nm to 1 565 nm. The L band extends approximately from 1 565 nm to 1 625 nm. The U band extends approximately from 1 625 nm to 1 675 nm. The most widely used band is the C band. The invention relates to fibers with chromatic dispersion compensation in the S band.\nIn the prior art, it is known to combine certain types of dispersion shifted optical fibers reducing crossed non-linear effects, known as non-zero dispersion shifted fiber (NZ-DSF), with dispersion compensating fiber (DCF), thereby achieving a transmission line with no dispersion over a wide spectral range.\nThe patent application EP 1219986 discloses using an optical fiber with chromatic dispersion compensation in the S band, but the examples described with two or three slices either have a ratio between the attenuation at 1 520 nm and at 1 500 nm greater than two or a chromatic dispersion to dispersion slope ratio that renders them unsuited to the object of the invention.\nThe invention proposes a chromatic dispersion compensating optical fiber which, either because of profiles with at least four slices or because of profiles with three slices having structures corresponding to particular ranges of index and radius, provides effective compensation in the S band of the chromatic dispersion of a certain range of line optical fibers, combined with attenuations at 1 500 nm and at 1 520 nm that are comparable. A particularly advantageous aspect is the very deeply buried nature of the index of the first buried slice (or of the only buried slice if there is only one) relative to the index of the central slice."}
{"text": "In superconducting magnets, variations in magnetic fields and temperature are common occurrences which often cause conduction faults in superconducting filaments. To compensate for such conducting faults, the magnetic windings are usually formed of wires in which the superconducting filaments are encased in a metal, such as silver, copper or aluminum, which is normally conductive across a wide range of temperatures, including superconducting temperatures. The metal electrically stabilizes the wire by shunting any portion of the filament(s) which has become non-superconducting. Furthermore, the encasing metal has good thermal conductivity, tending to transfer heat away from hot spots in the windings. Nevertheless, events, such as flux jumps, wire motion or eddy currents, in the superconducting filaments or encasing metal produce heat which may lead to a precipitous normalization of superconductive windings. Such temperature excursions occur during periods measured in microseconds, such short periods being insufficient for dissipation of the heat to the coolant, e.g., liquid helium.\nTo reduce the probability of normalizing temperature excursions, U.S. Pat. No. 4,171,464 proposes that particles or fibers of gadolinium oxide or gadolinium-aluminum oxide be incorporated in the metal, e.g., copper, that encases the superconducting filaments. These materials have high heat capacities at superconducting temperatures, e.g., below about 5.degree. K., and thus tend to absorb substantial amounts of locally produced heat. Furthermore, these materials undergo antiferromagnetic transitions within the superconducting temperature range, and when a decrease in the magnetic field occurs as a result of a conduction fault, the specific heat of the material increases, enhancing the adiabatic stabilization.\nHowever, adiabatic thermal stabilization through adiabatic demagnetization by compositions, such as the gadolinium and gadolinium-aluminum oxides which undergo antiferromagnetic transitions, is only realized if the effective magnetic field experienced by the material decreases. Thus, the material is less efficient in preventing precipitous normalizations in circumstances where a localized normalization of the superconducting filaments occurs but does not produce an immediate change in the magnetic fields to which the thermal stabilizer is subjected.\nFurthermore, because the magnetically affected contributions to heat capacity of materials, such as gadolinium oxides, are a substantial portion of the total heat capacity and because these components of their total heat capacity decrease with increasing magnetic field, these materials are less useful for thermal stabilization in superconducting applications where very intense magnetic fields are produced. As noted in the above-mentioned '464 patent, Gd.sub.2 O.sub.3 has substantially less heat capacity in a 2.4 Tesla (Tesla=10,000 Gauss) field than in a 0 Tesla field. Many present day superconducting magnet applications require substantially higher magnetic fields; for example, at the present time, superconductng magnets for fusion applications may produce magnetic fields in the range of 15 Tesla. At such magnetic field intensities, the heat capacities of materials, such as Gd.sub.2 O.sub.3, having large magnetic field-affected components of heat capacity, will be severely depressed.\nThere is a need for materials for use in superconducting applications which have high heat capacities and which are substantially unaffected by the magnetic field."}
{"text": "The present invention relates to an automatic sheet cutting apparatus which cuts off a sheet like a knit or a woven fabric into a predetermined shape. More particularly, the invention relates to a vacuum air leakage preventive device combined with an automatic sheet cutting apparatus so that evacuated air can be prevented from leaking through a clearance between a bottom surface of one of any brush-mounting bases aligned in parallel with a pair of endless chains and a top surface of the main body of the automatic sheet cutting apparatus.\nAny of those conventional chain-driven automatic sheet cutting apparatuses is provided with a number of brush-mounting bases which are aligned on a top surface of a sheet cutting table in parallel with a pair of endless chains installed on both sides of the main body of the automatic sheet cutting apparatus by way of being suspended on them. These brush-mounting bases sequentially move themselves relative to the rotation of a pair of endless chains. In the same way, a number of brushes secured on respective bases move themselves in conjunction with a cuttable sheet. Next, a control system stops the rotation of these endless chains at a predetermined position, and then activates a suction unit below the main body of the automatic sheet cutting apparatus to position the objective sheet on those brushes via a number brushes and those brush-mounting bases before permitting a cutting unit to cut the sheet on a cutting table surface into a predetermined shape by moving the cutting unit to the cutting position.\nOn the other hand, since there is substantial space between a pair of endless chains on both sides of the automatic sheet cutting apparatus, a supporting member is disposed on the main body of the automatic sheet cutting apparatus so that a number of brush-mounting bases aligned in parallel with and being suspended on these endless chains can slide themselves on the supporting member. This in turn results in an increased resistance against the sliding movement of those brush-mounting bases, and therefore, the apparatus needs to provide such a drive motor having substantial drive capacity.\nTherefore, in order to solve this problem, any conventional automatic sheet cutting apparatus is provided with sliding members or rollers below respective brush-mounting bases to shift those brush-mounting bases by causing those sliding members or rollers to slide or roll themselves on the supporting member.\nNevertheless, when these sliding members or rollers installed below those brush-mounting bases on the supporting member of the automatic sheet cutting apparatus slide or roll themselves on the supporting member, a narrow clearance is generated between the bottom surface of those brush-mounting bases and the supporting member. In consequence, when operating a suction unit to pneumatically evacuate a cuttable sheet on those brushes via those brush-mounting bases and brushes, evacuated air leaks out of this clearance, thus resulting in a lowered efficiency.\nTherefore, an object of the invention is to fully solve those technical problems inherent in any of those conventional automatic sheet cutting apparatuses by providing a novel system for preventing evacuated air from leaking out of a clearance between the bottom surface of each brush-mounting base aligned on a pair of endless chains and the top surface of the main body of the automatic sheet cutting apparatus."}
{"text": "Certain presentations and graphs require printing on a printable media, such as a paper sheet, of a dimension that is most beneficially presented in the form of a strip, for example, presentation charts used in project management. Project management is the planning and control of many activities that must be coordinated to achieve specific goals leading to the completion of an overall given objective. The project management process frequently uses a set of tools which incorporate charts and reports to detail the project for communication within the project team and with others.\nOne form of project management chart is referred to as a work breakdown structure (WBS) chart. A WBS chart is an organizational diagram type of chart depicting work packages comprising all of the principal elements of a project. Another form of project management chart, used for communication with the project team and with others, is a precedent network (Network) chart, sometimes erroneously referred to as a PERT chart. The Network chart displays activities required to produce the work packages depicted in the WBS chart and shows the relationships between the activities, i.e. the precedents and dependencies between the activities as they flow towards completion of a project.\nFor practical project management purposes, the WBS and Network charts are usually more than one standard sized page in width. To present either chart may require many standard sized pages to be joined one to the other. As the project evolves, the project management process results in the updating, modifying, and reproducing the project charts as a consequence of project progress. Thus, the WBS and network charts will change during the course of the project, consequently requiring the WBS and network charts to be produced frequently during the course of the project. To allow the charts to be printed on standard sized paper using conventional computer printers or photocopied onto standard sized paper, requires the joining of standard sized pages together to form a completed WBS or Network chart. Joining these pages is a time consuming process and usually requires input from the project manager to lay out the pages in a proper sequence prior to cutting and pasting them together to form the chart. Each time revised charts are produced during the life of a project, several copies are required each for key team members. The page assembly process to produce the charts is a frustrating task and often results in sloppy presentations even though the computer-generated data or the images printed or photocopied on each of the individual pages may be perfect."}
{"text": "Tertiary treatment of municipal sewage is a common wastewater application for ultrafiltration and microfiltration membranes; however, such systems need to be capable of operating on high suspended solids feed waters while having a long life with a minimum of chemical cleanings. Suspended solids that need to be removed may be materials that cause turbidity, such as bacteria, cysts and oocysts, viruses, colloidal material, such as iron oxides, clay, silt, sand and other insoluble impurities. Municipal sewage secondary treatment effluent typically has turbidity levels of 5 to 10 NTU with a suspended solids count of 10 to 20 parts per million (ppm). For membrane technology to be economically competitive in a tertiary treatment process, it should operate at sustained permeate flux rates of 15 to 30 gallons per square foot per day (gfd) with minimal energy consumption, while requiring chemical cleaning at a frequency of not more than once per month.\nPerhaps the biggest dilemma facing ultrafiltration and microfiltration membrane technology is fouling, where suspended solids and natural organic matter are adsorbed onto the membrane surface and result in a decrease in permeate flux. Decreased permeate flux rates result in higher overall capital and operating costs due to the lower output over a given period of time. The development of new membrane technologies has focused on reducing or eliminating the adverse affects of membrane fouling. A membrane element or cartridge that can produce high flux rates while minimizing fouling can have an extreme economical advantage over competing technologies.\nHistorically, such difficult applications as treating feed solutions high in organic and suspended solids have employed spiral wound, hollow fiber, capillary, or tubular element designs in a multitude of configurations. The most common of configurations is a submerged membrane system where the membranes, whether it be spiral wound or hollow fiber, are placed vertically in an open tank and filled with a feed solution to be filtered.\nVery generally, a spiral wound membrane cartridge contains a permeate carrier sheet, a membrane filter sheet that is adhesively bonded to the permeate carrier sheet (usually to both surfaces thereof to create an envelope about it), and a feed spacer sheet which separates two facing membrane filter layer sheets. Microfiltration (MF) and ultrafiltration (UF) membranes are typically formed of either polyethersulfone (PES), polysulfone (PSF), polyvinylidene fluoride (PVDF), or polyacrylonitrile (PAN) because these polymers are generally recognized in the industry to make membranes having high flux rates, good chemical resistance and good physical durability, which can be produced using conventional casting techniques. Other polymers such as polypropylene, polyethylene, and chlorinated polyethylene may also be used to construct UF and MF membranes. Spiral wound membrane cartridges exhibit good fouling resistance and offer economically attractive filtration with greater mechanical durability.\nA traditional submerged membrane system is shown in FIG. 6 of WO 00/78436 patent application (28 Dec. 2000) wherein a spirally wound membrane element is immersed in a tank that is filled with a body of water to be filtered and air is bubbled up through the membrane for the purpose of maintaining a clean membrane surface. The required transmembrane pressure (TMP) is supplied by a pump that creates a vacuum in addition to any contribution from the static head. Alternatively, static liquid heads alone have been used to generate feed pressures for submerged filtration, see U.S. Pat. No. 5,916,441. These systems will typically also employ support frames and manifolds. Materials of construction are typically limited to 316 stainless steel for these type of systems, which is a significant capital expenditure. The use of 316 stainless steel also limits the amount of exposure to certain corrosive water sources, such as cleaning solutions and aggressive feedwaters such as seawater.\nFactors which contribute to membrane fouling include water chemistry, suspended solids concentration (TSS), membrane chemistry and transmembrane pressure (TMP), are used to combat same. Membrane chemistries vary between polysulfone (PSF), polyethersulfone (PES), and polyvinylidene fluoride (PVDF); however, all are generally similar in their fouling tendencies. The more hydrophilic a membrane chemistry is, the more resistant it is to fouling, as fouling constituents are generally hydrophobic.\nMembrane fouling typically occurs through two different mechanisms; pore impregnation and cake-layer formation. Pore impregnation occurs when organic and suspended solids penetrate and impregnate the pores of a membrane, thus narrowing the effective diameter of the pores and increasing the required net drive pressure for permeate production. Cake-layer formation involves large colloidal compounds and suspended solids accumulating on the membrane surface and forming a discrete layer of organic and suspended solids. Both fouling mechanisms result in decreased permeate flux.\nCommon methods to offset the effects of fouling for submerged membrane systems include air bubbling, backwashing, and chemical cleaning. Air bubbling, whether it be for a spiral wound or hollow fiber membrane, is quite effective at removing particulate matter from the membrane surface; however, it is an energy intensive process that often comprises 50% or more of the total energy costs associated with a submerged membrane system. The economical advantages of a membrane system that can eliminate or greatly reduce the use of air bubbling without increasing chemical cleaning frequency are obvious.\nWhen the effects of fouling become irreversible, membrane cartridges require extensive cleaning with a variety of chemical cleaners. Existing submerged membrane systems often rely on extensive static soaking with an oxidant-based, caustic-based, and/or acid-based cleaner. Air bubbling is discouraged during cleaning and is often discontinued during the cleaning process so as to avoid excessive foaming and/or release of harmful vapors to the atmosphere. However, static cleaning has been found to be much less effective than dynamic cleaning wherein the added shear force enhances the removal of organic and suspended material from membrane surfaces.\nAnother problem associated with submerged membrane systems is the accumulation of solids in the process tank and the inefficient removal thereof. With submerged membrane systems, the suspended solids level inside the tank increases over time, due to accumulation, until a steady state value is reached. This steady state suspended solids level is directly proportional to the recovery of which the system operates at. For example, a system that operates at 90% recovery will have a steady state suspended solids value about ten times greater than that of the influent level. Since membrane fouling is directly proportional to suspended solids concentration, it is obvious that membrane fouling will be greater when the suspended solids level has reached its steady state concentration value; however, it is highly undesirable to lower the system percentage recovery for the sake of reducing fouling as such would drastically lower the overall efficiency of the system, potentially rendering it impractical and economically unattractive. Moreover, after every backwash, solids removed from the membrane surface are mixed directly back into the bulk tank solution, further complicating the solids build-up problem. Because large, submerged membrane systems utilize process tanks having very large holding volumes, it is impractical to drain said tanks on a frequent basis in order to purge the accumulated solids from the system.\nWith the foregoing in mind, work on improvements was undertaken."}
{"text": "Wireless scanners, or readers, may detect the presence of tags. If the scanners fail, replacing the scanners may be cumbersome, resulting in some tags going undetected."}
{"text": "Because silicon nitride sintered products are excellent in various properties such as mechanical strength, heat resistance, etc., they are used as high temperature structural materials, such as for parts of gas turbines. However, on the other hand, silicon nitride has a poor sinterability because of having a highly covalent bonding nature, and, consequently, it is difficult to obtain sintered products having a high density and high strength, employing silicon nitride per se. Therefore, in the case of sintering silicon nitride, sintering aids including MgO, Al.sub.2 O.sub.3, oxide of rare earth elements, etc., have been used hitherto. However, sintered products using MgO as a sintering aid have inferior oxidation resistance because MgSiO.sub.3, etc., are formed in the oxidizing atmosphere. In sintered products using Y.sub.2 O.sub.3 as a sintering aid, the oxidation resistance deteriorates extremely if YSiO.sub.2 N and Si.sub.3 Y.sub.2 O.sub.3 N.sub.4 are formed as the second phase (the meaning of which is hereinafter explained) as opposed to forming Si.sub.2 ON.sub.2, Y.sub.2 Si.sub.2 O.sub.7, or Y.sub.5 Si.sub.3 O.sub.12 N. Formation of the second phase is affected by the amount of Si.sub.3 N.sub.4, SiO.sub.2, or Y.sub.2 O.sub.3, the kind and the amount of other sintering aids, the sintering atmosphere, etc., and it is difficult to correctly control these factors. In the case of using CeO.sub.2 as a sintering aid, the second phase composition having good oxidation resistance is not formed in an Si.sub.3 N.sub.4 -CeO.sub.2 system, because CeO.sub.2 is reduced to Ce.sub.2 O.sub.3 during sintering and the produced oxygen reacts with Si.sub.3 N.sub.4 to form SiO.sub.2. In the case of using other oxides of rare earth elements as sintering aids, compositions having improved oxidation resistance have not yet been obtained."}
{"text": "The primary source for oxygen for a grill comes from below the burner. Conventional grills have vents located near the bottom of the grill, below the burner, to allow for air flow which supplies the necessary oxygen. However, if the airflow through these vents is at too high a velocity, the flame at the burner can be blown out. This can cause conventional grills to be virtually useless. in high wind environments.\nAdditionally, other areas of high velocity airflow can also disrupt the burner operation. One such area is the venturi assembly. Conventional grills have a venturi tube for mixing propane and air which is generally located outside of the grill tub. High wind conditions can result in this mixture being diluted too much, resulting in the burner flame going out. Additionally, high velocity airflow in the area directly above the burner can disrupt the burner operation. This can occur either when the grill hood is open, allowing the wind to access the top of the burner through the heat diffusing material, which typically consists of rocks or ceramic bricketts. Wind can also affect burner operation through vents in the grill hood when the hood is closed. In any of these cases, the high wind conditions can cause the flame at the burner to go out, and thus the grill must be re-lit, if possible, in order to continue operation.\nBased on the foregoing, there is a need for a grill which overcomes these problems and protects the flame at the burner from high velocity airflow."}
{"text": "Electromagnetic devices such as motors, generators, magnets, and solenoids, use stacked laminations made from silicon-iron (Si--Fe) for example for providing efficient flux paths for the operation thereof. The laminations must be suitable stacked together either using mechanical fasteners or by welding. Exemplary mechanical fasteners include straps, keys, and rivets suitable positioned in the laminations for holding them together. However, the laminations are typically pre-machined for accepting the mechanical fasteners such as by being punched or providing apertures therefor, which premachining induces stress in the laminations. A pulsed annealing operation is therefore typically required for reducing the residual stresses for minimizing magnetic core losses therefrom. The location of the fasteners themselves may also lead to inefficient flux paths within the laminations decreasing the efficiency of the device.\nWelding, on the other hand, necessarily heats the laminations which adversely affects the magnetic properties thereof and may distort the laminations. This may require a pulsed-heat treatment to reclaim the lost magnetic properties and remove the distortions.\nConventional TIG welding is one process known to provide excessive heat into the laminations necessarily requiring the pulsed-heat treatment thereof. Excessive heat input into the laminations may also cause physical distortion thereof which must be suitably corrected in order to obtain a satisfactory product.\nIt is also known to use a conventional CO.sub.2 laser for welding stacked laminations for reducing the heat input thereto and reducing both thermal distortion thereof and the degradation of magnetic properties therein. However, CO.sub.2 lasers when conventionally used for welding are typically operated as continuous-wave lasers and therefore continuously input heat into the welded member although such heat input is significantly less than that from conventional TIG welding. Continuously applied welding heat and the resulting continuous, substantially uniform weld bead may still degrade the magnetic properties of the laminations from excessive heat."}
{"text": "As taught in U.S. Pat. No. 3,489,687 (Inamorato et al.), U.S. Pat. No. 4,276,205 (Ferry), and U.S. Pat. No. 4,659,565 (Smith et al.), it is known that mixed tert-amine oxides containing alkyl groups of different chain lengths can be used in conjunction with other materials, such as other surfactants, detergent builders, and/or other additives, to prepare liquid or solid detergents.\nThe mixed tert-amine oxides which have been synthesized by conventional techniques, i.e., the typical aqueous solutions of mixed tert-amine oxides, present no particular problems in the preparation of liquid detergents in which their water content is not a liability. However, they have to be dried in order to be used in the preparation of solid detergents.\nThe technique used by Inamorato et al. to dry their aqueous tert-amine oxides is to mix them with sufficient amounts of hydratable inorganic salts, such as sodium sulfate or sodium tripolyphosphate, to react with substantially all of the water in the solution. Using this technique, they form a friable amine oxide/hydrated salt mixture which they can granulate to obtain a free-flowing detergent composition.\nWhen Ferry uses tert-amine oxides to prepare solid detergents, he spray-dries the aqueous solutions--a technique that is effective in removing the water but has the unfortunate side-effect of subjecting the amine oxides to temperatures at which they are apt to decompose and form amines that are corrosive to the skin."}
{"text": "Semiconductor devices are becoming smaller and more dense with the evolution of new technology. However, increases in circuit density produce a corresponding increase in overall chip failure rates at a time when chip failure rates must decrease to remain competitive. Chip manufacturers are therefore challenged to improve the quality of their products by identifying and eliminating defects which produce defective chips known as fails. Whereas significant improvements are being made to eliminate systematic defects by reducing process variability. Process improvements alone are not sufficient to eliminate all the random defects which effect both yield and reliability. Historically, screening techniques have been employed to improve product failure rates to acceptable levels by culling out many of these random defects.\nSingle chip or multi-chip packaging strategies require the ability to join or replace individual chips from an assembled chip module. This should of course be done without altering the integrity of the remaining chips or capacitors or other electronic components. Chip changing is also required for yield management, upgrade capability, and engineering changes, to name a few. Chip replacement typically involves three distinct steps: chip removal, chip site solder dress, and new chip rejoin. Different tools and processes are used in the chip removal or chip join process.\nU.S. Pat. No. 3,283,124 (Kawecki), discloses a selective heating apparatus where energy from an infrared source is focused through the use of a special focusing means.\nU.S. Pat. No. 4,160,893 (Meyen et al.) discloses an individual chip joining machine, where an IR lamp is used to heat the upper surface of a chip through a flat quartz plate.\nU.S. Pat. No. 4,278,867 (Tan) discloses a system for chip joining by short wavelength radiation. Tan uses a laser to heat the upper or back-side surface of a chip during the process.\nIBM Technical Disclosure Bulletin, \"Luminous Dual Reflection Lamp\" Vol. 20, No. 6, Pages 2141-2142 (November 1977), discloses a lighting system for uniformly scattering light without projecting either characteristic central hot spots or concentrated filament images. This is achieved by using a flat clear quartz plate, having a reflecting convex portion and a concave reflector.\nPuttlitz, et al., \"Replacing Flip Chips on MLC Multichip Modules Using Focused Infrared (IR): The Basics\", ISHM '92 Proceedings, Pages 384-390 (October 1992), discloses a means of locally heating a target chip, site dressing the chip site, and placing and reflowing a replacement chip. This is done by using an infrared (IR) quartz-halogen lamp within a metal reflector housing. The IR lamp is typically located at or near the parabolic reflector focal point so that its energy can be directed and focused onto the backside of the target chip. Heating can, however, be confined to the target chip by introducing a substrate shield that prevents irradiating the surrounding MCM surface and attachments. This assures that only the target chip is irradiated (i.e., heated) when the IR lamp is energized."}
{"text": "Hot beverage makers, such as coffee makers, have been known and sold for many years using various brewing techniques. The typical and traditional coffee maker includes a stand or tower that has a warming plate forming the bottom or base of the tower with a filter basket located above the warming plate. The interior of the tower defines, at least in part, a fresh water reservoir. Such coffee makers further include a fluid reservoir, such as a glass carafe, that rests on the warming plate beneath the filter basket. Alternatively, the fluid reservoir may be an insulated carafe, in which case the warming plate is typically omitted.\nIn use, an operator fills the carafe in order to transfer water to the fresh water reservoir. The water is heated and passed through the filter basket, which includes the grounds to be infused. The brewed beverage then flows from the basket into the carafe. The beverage is maintained at an elevated temperature via the warming plate upon which the carafe rests, in the case of a glass carafe, or by the insulating properties of the carafe in the case of an insulated carafe.\nA new variation of coffee maker has been developed wherein a brewed beverage tank is included, such as the coffee makers described in commonly assigned U.S. Pat. No. 6,564,975 to Garman, issued May 20, 2003, and U.S. Pat. No. 6,681,960 to Garman, issued Jan. 27, 2004, the contents of which are incorporated herein in their entirety. Briefly, the brewed beverage tank holds the filter basket above a reservoir portion. Hot water passes through the filter basket and a material to be infused (e.g., ground coffee beans). The brewed beverage is then collected and held in the reservoir portion of the brewed beverage tank (“brew tank”). A dispenser actuator is depressed that opens an outlet port in the reservoir. A user simply actuates the dispenser actuator with a mug or cup and the brewed beverage passes through the outlet to the operator's container.\nIn order to allow a user to determine the amount of brewed coffee within the device, a transparent window on the coffee holding reservoir can be used. A transparent vertical slot-type window on the side of the holding reservoir can also be used. Also known is the use of a buoyant ball contained within a transparent column, the transparent column containing coffee at the same level as in the holding reservoir and the ball floating on the surface of the coffee in the column. The transparent column typically includes markings that correspond to the quantity of coffee in the reservoir, such that the quantity of coffee in the reservoir may be determined by reading the marking nearest the floating ball.\nHowever, these transparent windows, slots, and columns give a somewhat rough estimate of how much coffee remains and often do not give a clear indication of the exact amount. In addition, the amount of coffee remaining in the coffee maker is often difficult to gauge as the coffee may be difficult to see against the background of the coffee maker which is often somewhat similar in color to the coffee itself. Further, the glass or plastic of the these transparent windows, slots, and columns may be quickly stained by the coffee, such that the transparency is greatly reduced and the amount of coffee becomes difficult or impossible to determine. Because of the above problems, even when the amount of coffee may be determined at a short distance from the coffee maker via these transparent windows, slots, and columns, it may be difficult or impossible to determine the amount of coffee at a greater distance, such as across a room.\nThere is a need, therefore, for a coffee level indicator that is accurate, reliable, clearly indicates the amount of brewed coffee remaining in the holding reservoir, is unaffected by staining, and is readable at a distance. The operation and structure of a brew tank with an integrated fluid gauge in accordance with the present invention would solve one or more of these or other needs."}
{"text": "Traditionally prescribing ophthalmic lens to a wearer involves fitting a potential wearer having an eyesight defect with ophthalmic lenses having different optical characteristics and then selecting the lenses which provide the best optical performance for the wearer. Such techniques suffer the drawback that the wearer has to go through the process of putting off and taking off a number of spectacles and may only be exposed to a limited selection of lenses which may not include the optical lens optimal for correcting the particular eyesight defect.\nTo address such issues various solutions based on virtual reality techniques which simulate the optical effect of an ophthalmic lens have been proposed. For example, WO 2007/056795 describes a method of simulating an optical defect of a selected ophthalmic lens design which includes retrieving simulation data for the selected ophthalmic lens design and processing the simulation data to generate a virtual output image simulating the optical defect.\nA drawback of the system described in WO 2007/056795 is that it requires complex data processing techniques and a high calculating power to provide an average quality output image. This leads to prohibitive costs and complex system requirements for implementation of the virtual simulation method."}
{"text": "This invention relates in general to the construction of sewing machines and, in particular, to a new and useful multineedle clamps having a plurality of bores, the number of which corresponds to the number of needles and which are provided side-by-side and parallel to each other between the lateral walls of the needle guard body and intended to receive a needle butt."}
{"text": "(a) Field of the Invention\nThe present invention relates to an organic light emitting display and a manufacturing method thereof.\n(b) Description of Related Art\nAn organic light emitting display (OLED) is a self emissive display device, which displays images by exciting an emissive organic material to emit light. The OLED includes an anode (hole injection electrode), a cathode (electron injection electrode), and an organic light emission layer interposed therebetween. When the holes and the electrons are injected into the light emission layer, they are recombined and pair annihilated with emitting light. The light emission layer further includes an electron transport layer (ETL) and a hole transport layer (HTL) as well as an electron injecting layer (EIL) and a hole injecting layer (HIL) for enhancing the light emission. Each pixel of the OLED includes two TFTs, i.e., a switching TFT and a driving TFT. The current for light emission is driven by the driving TFT and the mount of the current driven the driving TFT is controlled by the data signals from the switching TFT.\nA plurality of pixels of the OLED, each including an anode, a cathode, and a light emission layer, are arranged in a matrix and driven in passive matrix (or simple matrix) addressing or active matrix addressing.\nThe passive matrix type OLED includes a plurality of anode lines, a plurality of cathode lines intersecting the anode lines, and a plurality of pixels, each including a light emission layer. The selection of one of the anode lines and one of the cathode lines cause light emission of a pixel located at the intersection of the selected signal lines.\nThe active matrix type OLED includes a plurality of pixels, each including a switching transistor, a driving transistor, and a storage capacitor as well as an anode, a cathode, and a light emission layer. The OLED further includes a plurality of gate lines transmitting gate signals and a plurality of data lines transmitting data voltages. The switching transistor is connected to one of the gate lines and one of the data lines and transmits the data voltage from the data line in response to the gate signal. The driving transistor receives the data voltage from the switching transistor and drives a current having a magnitude determined depending on the difference between the data voltage and a predetermined voltage such as a supply voltage. The current from the driving transistor enters the light emission layer to cause light emission having an intensity depending on the current. The storage capacitor is connected between the data voltage and the supply voltage to maintain their voltage difference. The gray scaling of the active matrix type OLED is accomplished by controlling the data voltages to adjust the current driven by the driving transistor. The color representation of the OLED is obtained by providing red, green and blue light emission layers.\nIn the meantime, the OLED is classified into top emission type and bottom emission type depending on the light emitting direction. The top emission type OLED includes a transparent cathode usually made of indium tin oxide (ITO) or indium zinc oxide (IZO) and an opaque anode, while the bottom emission type OLED includes an opaque cathode and a transparent anode. The relative positions of the anode and the cathode can be altered if required.\nSince ITO and IZO have high resistivity, the top emission type OLED may add an auxiliary electrode having low resistivity to the cathode. However, the addition of the auxiliary electrode requires an additional lithography step that may complicate the manufacturing process and increase the manufacturing cost."}
{"text": "This invention relates to high-intensity-discharge (HID) lamps and, more particularly, to improved electron emissive material for the electrodes of such lamps.\nIn U.S. Pat. No. 3,708,710 dated Jan. 2, 1973 is disclosed a high-intensity-discharge sodium-mercury vapor lamp which utilizes dibarium calcium tungstate as electron emissive material. Such material has been used in so-called dispenser cathodes and U.S. Pat. No. 3,434,812 dated Mar. 25, 1969 discloses the use of dibarium calcium tungstate or dibarium strontium tungstate as an emissive material in a dispenser cathode.\nDibarium calcium molybdate is known for use as a getter layer material in conjunction with an incandescent lamp, as disclosed in U.S. Pat. No. 3,266,861, dated Aug. 16, 1966. In addition, high-pressure sodium-mercury vapor lamps have in the past utilized as electron emissive material a mixture of several oxide phases comprising thorium dioxide, barium thorate, dibarium calcium tungstate and barium oxide. This mixture of oxide phases is quite sensitive to the atmospheric contaminants with the result that even a brief exposure to the air can result in a relatively large pickup of water and carbon dioxide by the emission mixture, which contaminants are rather difficult to remove. In such a mixture, the thorium dioxide serves as a matrix for the more active oxide emitters such as the barium oxide, dibarium calcium tungstate and barium thorate."}
{"text": "Applications and processes executing in enterprise cloud based environments may be hindered because of the lack of flexible resource scaling and provisioning of resources. Applications and processes executing in the cloud may have to adapt to varying resource requirements based on customer and business needs. Provisioning computing storage resources at the time usage demand increases may be inefficient because it may take significant time and resources to allocate and initialize the storage which may cause service interruptions. Thus, operators may be faced with the choice of expensive over-provisioning of computers resources, or reactively provisioning computer resources when demand increases occur which may result in poor service response and unplanned downtime. Further, some resources, such as virtual machine data storage, may have provisioning characteristics where reactive provisioning may cause significant negative impact on system performance. Thus, it is with respect to these considerations and others that the invention has been made."}
{"text": "E-commerce has grown exponentially over the past decade. In addition, consumers are increasingly using the Internet to perform various other transactions such as banking, information gathering, etc. However, security associated with these on-line transactions has always been a cause for concern for all users. These concerns about security has both hindered growth of on-line transactions and undermined consumer confidence in on-line transactions. One of the major areas of concern for consumers is protecting their sensitive information from getting into the hands of criminals.\nOne of the predominant ways of stealing sensitive consumer information such as credit card numbers, birth dates, social security numbers, etc. is to hack into a merchant website or into a consumer's computer and steal this information. It is more difficult to compromise a merchant website since most of the reputed merchants use sophisticated technology to keep out the criminals. However, consumer communication devices, e.g., a computer, often lack adequate protection against unauthorized access making them an easy target for stealing sensitive information.\nHackers use numerous techniques such as, Cross-Site Scripting (CSS), Cross-Site Request Forgery (CSRF), Phishing, Cross-frame Scripting, Man-in-the-Middle (MIM), browser cache attacks, cookie exposure, etc. to gain access to sensitive consumer and/or merchant information. One of the predominant methods to conduct these attacks is via an Internet browser, e.g., Internet Explorer, that may be resident on a consumer computer this is usually used to conduct on-line communications and e-commerce transactions.\nConventional techniques for preventing such attacks include providing a sand-boxed environment for conducting on-line transactions. However, even the sandbox solutions are not immune from CSRF attacks that exploit the trust a website has for a user of the website.\nWhat is needed is a system and method to alleviate and mitigate these threats and to provide a secure environment for conducting on-line transactions."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate in general to compositions and methods that are useful for incorporating certain alloying compounds into metal substrates. More particularly, embodiments of the present invention relate to fluid compositions comprising certain alloying compounds which may be sprayed onto metal substrates and then heated, for example by a laser, to temperatures sufficient to create a molten mixture and to disperse the alloying compound into the metal substrate.\n2. Description of Related Art\nCompositions do exist for bonding and/or cladding the surface of metal substrates with transition elements by using high intensity lasers. The interface between the cladding and the surface of the metal substrate is clearly defined with no substantial intermixing between the two sections. These compositions or the methods by which they are applied do not alter the chemical or physical composition of the metal substrate itself to improve characteristics.\nSuch compositions are also unfortunately generally in the form of dry metal powders and their application to metal substrates is often difficult and wasteful of the dry powder. Also, the application of the dry powder directly to the surface of the metal substrate involves several laborious steps to pretreat the metal substrate prior to application of the metal powder to ensure that it will adhere to the metal substrate. In addition, given the difficulties in working with dry powders, dry powder deposition techniques are usually limited to applications involving small surface areas.\nAccordingly, there is a need in the art to provide compositions and methods for incorporating alloying elements into metal substrates and also in a manner which overcome the drawbacks of existing compositions and methods of cladding metal substrates, while providing for economical and easy application."}
{"text": "The following description of the background of the invention is provided simply as an aid in understanding the invention and is not admitted to describe or constitute prior art to the invention.\nThe present invention relates generally to the field of audience interaction systems. Specifically, the present invention relates to audience messaging systems.\nMethods for sending electronic messages such as the short message service (SMS) or text messaging allows a user to send simple text messages from a mobile device, computer or handheld device to a recipient device such as a mobile phone. FIG. 1 illustrates a general communications architecture for mobile device to mobile device messaging. When a text message is sent from a mobile device 10 it is first received by a mobile switching center (MSC) 20. The MSC 20 then forwards the message to a short message service center (SMSC) 30. The SMSC 30 routes SMS messages within a wireless provider's network. The text message is saved at the SMSC and then forwarded to a second MSC 40 associated with the recipient device 50 which then routes the text message to a recipients mobile device 50. SMS messages are stored by the SMSC until the destination device is available to receive the message. Thus, the save and forward capability of the SMSC guarantees deliver of the SMS message.\nText messaging is an increasingly popular communication method. Several commercial applications for text messaging exist. For example, companies use broadcast text messages to send information to multiple consumers. Text messaging is used by television programs to allow viewers to participate in interactive polls. In addition, venues for concerts and sporting events use displays such as giant screens, “JUMBOTRONS®, “ribbon boards”, video displays in the venue, electronic marquees and scoreboards to display text messages sent from audience members. FIG. 2 shows an exemplary communication architecture for a conventional venue messaging system. A user sends a message on her mobile device 10 to a number associated with the venue's display 80. The message is transmitted to a MSC 20 which routs the message to a SMSC 30. The SMSC 30 stores the message and then routes the message to a message server 60. The message server 60 passes the message through a content filter 70. The filtered message is then shown on the venue display 80.\nIt is common for multiple events such as concerts or sporting events to occur at different venues at the same time. Thus, a system and method is desired that will allow an audience member located in one venue to send a message to be displayed to audience members attending a different event at a separate venue."}
{"text": "The present disclosure relates to a light emitting device and a method of manufacturing the light emitting device.\nFor a light emitting device, for example, Japanese Unexamined Patent Application Publication Nos. 2013-219397 and 2016-219743 each disclose a structure which includes a cover member disposed in the surrounding of a light emitting element."}
{"text": "Semiconductor lasers are commonly used for low power reading operations in optical data storage systems, such as CD-ROM, audio compact disk and video disk systems. For these applications, red (620-690 nm) AlGaInP/GaAs laser diodes are favored over near infrared (780-880 nm) AlGaAs/GaAs laser diodes, because of their shorter wavelength and hence better resolution, although the latter are also sometimes used because of their lower unit price. In addition to such wavelength considerations, laser diode application often favor one particular polarization, either transverse electric (TE) or transverse magnetic (TM), over the other. Moreover, TE mode lasers typically have higher gain, higher output power and higher operating temperatures than TM mode lasers, which may be important factors in certain applications. In light of all of this, it would be advantageous to be able to selectively tailor the semiconductor laser materials to not only produce a particular emission wavelength but also to favor a preferred polarization at that selected wavelength. Unfortunately, for the AlGaInP/GaAs material system it has been found that prior laser diodes of this type operate in the TE polarization mode only for the longer 650-690 nm emission wavelengths and that only the TM polarization mode is available for the shorter 620-650 nm wavelengths. To see why this should be the case, it will be useful to consider the effects of material composition on bandgap energy and strain, and of strain on the emission polarization.\nFIG. 1 shows the bandgap energies and corresponding wavelengths for unstrained In.sub.x Ga.sub.1-x P material in the range from x=0.2 to x=0.8. (The wavelengths are derived from the relationship E.sub.g =hc/.lambda., where E.sub.g is the bandgap energy, h is Planck's constant, c is the speed of light in vacuo, and .lambda. is the wavelength, and thus assumes that substantially all of emitted photon energy is contributed by the bandgap energy.) It can be seen from FIG. 1 that the alloy composition In.sub.0.5 Ga.sub.0.5 P has a bandgap of about 1.9 eV and emits light at about 650 nm. In order to produce laser light emission with a longer wavelength than 650 nm, one can decrease the bandgap of the active lasing material by changing its composition, for example by increasing the proportion of indium. Thus, In.sub.0.58 Ga.sub.0.42 P has a bandgap of about 1.8 eV and a corresponding emission wavelength of about 690 nm. Likewise, in order to produce a wavelength shorter than 650 nm, one can change the composition of the material to increase its bandgap. While this could be done by adding aluminum to the composition, as is done with quantum barriers and cladding layers in the laser diode heterostructure, in order to provide sufficient carrier confinement for room temperature CW operation, the bandgap increase for the quantum wells is normally obtained by decreasing the proportion of indium in the composition. For example, In.sub.0.42 Ga.sub.0.58 P has a bandgap of about 2.0 eV and a corresponding wavelength of about 620 nm.\nFIG. 2 shows the amount of strain for various compositions of In.sub.x Ga.sub.1-x P layers on a GaAs substrate. The material is lattice-matched to the substrate for x.apprxeq.0.48. Increasing the proportion of indium results in a larger crystal lattice so that the material is under compressive strain. For example, when x.apprxeq.0.62, the alloy composition is under about 1.0% compressive strain. Likewise, decreasing the proportion of indium produces a smaller lattice size and tensile strain. Thus, when x.apprxeq.0.35, the composition experiences -1.0% compressive strain (i.e., +1.0% tensile strain). Replacing part of the gallium with aluminum has only a negligible effect on strain. For useful lifetimes, a maximum strain in the active region's quantum wells and barriers of about 2% (i.e., compositions where 0.2.ltoreq.x.ltoreq.0.8) is a practical limit for most devices. The critical limit for strain depends on the thickness of the strained layers. Thus, the cladding layers and other relatively thick layers outside of the active region require significantly less strain and may need to be substantially lattice-matched to the substrate to avoid dislocations and other lattice defects. One can see from a comparison of FIGS. 1 and 2, that compositions providing bandgaps less than 1.9 eV and corresponding emission wavelengths longer than 650 nm are under compressive strain, while compositions providing bandgaps greater than about 1.9 eV and corresponding wavelengths shorter than 650 nm are under tensile strain.\nFor GaInP lattice-matched to a GaAs substrate, the band-structure has \"light\" hole and \"heavy\" hole valence bands that are degenerate at k=0 (where k is the wavenumber of carriers in the lattice). There is an energy gap E.sub.g between these valence bands and the conduction band. TE polarized light is emitted when an excited electron in the conduction band recombines with a hole in the heavy hole valence band, whereas TM polarized light is emitted when an electron combines with a light hole. Because the differential gain of the TE mode is larger than that of the TM mode, a laser mode with a lattice-matched InGaP active region will generally lase in the TE mode. An In.sub.x Ga.sub.1-x P layer with x&lt;0.48 will be under tensile strain and will emit TM polarized light. However, a compressively strained layer (x&gt;0.48) will emit TE polarized light due to the dominance of the heavy hole band. Prior laser diodes emitting at the shorter 620-650 nm wavelengths have tensile strained quantum wells and thus operate in the TM mode, while those emitting at the longer 650-690 nm wavelengths are compressively strained and operate in the TE mode.\nWhen reading stored data from optical disks, a very low random intensity noise (RIN) level from the laser illumination source is needed for a good signal-to-noise ratio at the detector. However, in such data storage systems, external reflective feedback typically occurs, in which some fraction of the emitted laser light is coupled back into the laser diode. The phase of this feedback varies in time due to vibrations in the system, other slight, but essentially random, movements of various system components, and the reading of phase-encoded data stored on some types of disks. These phase variations in the feedback cause modal instabilities in the laser cavity, resulting in random fluctuations (i.e., noise) in the laser's output power.\nAlthough this feedback sensitivity of the laser is also dependent to some extent upon the amplitude of the optical feedback, such that the noise in the emitted output of the laser could be reduced by using an optical design that minimizes the amount of reflected light coupling back into the laser, such a design normally requires completely separate illumination and collection optics, thereby leading to a much larger and heavier (and hence mechanically slower) read head containing such optics.\nAnother way to reduce the feedback sensitivity of a laser diode is to modulate the laser at high speed. The high speed modulation renders the laser insensitive to feedback by destroying phase coherence between the reflected beam and the oscillating modes of the laser. Typically, the modulation rate needed for this purpose is on the order of several hundred megahertz. Unfortunately, the additional circuitry needed to bias and modulate the laser diode at high frequency considerably complicates the overall system and makes it much more expensive.\nA better solution is to provide a laser diode that inherently self-pulsates, emitting a high frequency stream of successive optical pulses when driven by a DC (unmodulated) bias. Self-pulsating laser diodes have been demonstrated, as illustrated for example in U.S. Pat. Nos. 4,961,197 (Tanaka et al.); 5,003,549 (Mitsui et al.); 5,111,469 (Narui et al.); and 5,416,790 (Yodoshi et al.). Most of these prior self-pulsating lasers are of the near-infrared AlGaAs-type. Those of the AlGaInP-type are characterized by the same wavelength-polarization constraints imposed by strain as their non-pulsating counterparts.\nThe self-pulsating laser diodes usually achieve their rapid pulsation operation by incorporating a saturable absorbing layer into the laser structure. This absorbing layer is positioned outside of the active region, but close enough to the active region that it overlaps the propagating optical mode. The active region itself can be a single active layer (typ., .apprxeq.50 to 150 nm thick), a single quantum well (SQW) structure (typ. well .apprxeq.10 to 30 nm thick), a multiple quantum well (MQW) structure (typ. 2 to 4 wells, each .apprxeq.5 to 20 nm thick), or a superlattice structure (effective thickness .apprxeq.15 to 30 nm). The saturable absorbing layer in the neighborhood of the active region is composed of a material (such as GaInP) and has a thickness (e.g., it may form a 10 to 40 nm thick quantum well) having an effective bandgap energy that is smaller than (or at most substantially equal to) the energy corresponding to the lasing wavelength, thereby giving the absorbing layer a high intrinsic absorption coefficient at the lasing wavelength. Accordingly, when the laser diode is first energized, the absorbing layer initially absorbs the light generated in the active region quite strongly, so that the photon density in the active region is kept at a level below the lasing threshold and the density of excited electron-hole pairs in the active region is allowed to build. Carriers generated in the absorbing layer due to the absorption of light from the active region are confined there by the higher bandgap energy of adjacent cladding layer material (such as AlInP) both above and below the absorbing layer. The resulting collection of carriers in the absorbing layer causes a drop in its absorption coefficient as the absorbing layer saturates, thereby reducing the loss in the laser cavity and allowing the sudden onset of strong lasing as the photon density in the active region rises above the threshold level. The intense lasing mode rapidly depletes the carriers in the active region, so the device quickly stops lasing again. The cladding material above the absorbing layer typically has a ridge structure with lower bandgap material (such as GaAs) regrown on each side of the central ridge. Carriers confined in the absorbing layer by the cladding layer below and the ridge above are still free to diffuse laterally within the absorbing layer until they reach the areas on each side of the ridge structure, where they can fall out of the absorbing layer into the surrounding lower bandgap material. The absorbing layer can thus be quickly emptied of the accumulated carriers by means of diffusion once lasing has stopped, so as to restore the absorption coefficient to its original peak value and again become a strong absorber of light generated by the active region. Hence, the cycle starts all over again. The modulation rate attained in such self-pulsating devices is typically several hundred megahertz, but can be as high as several gigahertz, which is more than enough to reduce the laser's feedback sensitivity and obtain very low RIN emission.\nIt is an object of the invention to provide AlGaInP/GaAs laser diodes that operate in the TE polarization mode for 620-650 nm emission wavelengths.\nIt is another object of the present invention to provide such TE mode laser diodes with self-pulsating operation."}
{"text": "This invention relates to battery chargers.\nLow voltage rechargeable batteries are used in many appliances, for example electric dry shavers. The appliance may include at least one rechargeable battery and an integral charger comprising a charging circuit including AC input terminals, a transformer, and a full wave current rectifier. A pilot lamp is desirably connected for energization during charging to indicate normal operation of the charging circuit. Extinction of the lamp signals some malfunction of the charging operation, such as lack of mains voltage or poor contacts.\nIn a compact hand-held appliance, such as an electric dry shaver, space and weight are strictly limited. Consequently the batteries, transformer, current rectifier and even the charging terminals for connection to the AC mains are designed to minimum dimensions and weight. The power output of the transformer is thus very limited and must be used with maximum efficiency to permit recharging to be accomplished within an acceptable short time period."}
{"text": "A wig has been manufactured in such a manner that a hair segment to be transplanted is folded in two, which is one by one transplanted onto a three-dimensional thick base by handwork. When one folded hair segment is transplanted on the base, it looks as if two hairs are transplanted. Several folded hair segments may be transplanted at one time.\nHowever, such a prior art handwork is extremely inefficient. It would take two or three weeks or more to transplant 20,000 hairs, for example. When a wig is manufactured in foreign countries in order to save labor costs, it tends to increase inferior products and reduce a production yield.\nSome attempts have been made to develop automated wig manufacturing systems, but no success has been achieved.\nAccording to the study by the present inventor, the greatest difficulty in automation of wig manufacturing exists in providing high-precision transplanting pitch. Human hairs are, in their natural condition, spaced from each other by less than 1 mm, or normally of the order of 0.5 mm, so that it will be desired to determine the transplanting pitch as such.\nHowever, it is quite difficult, like a divine work, to operate an extremely thin needle at a pitch of lower than 1 mm along a predetermined line. When the needle should be wobbling even a little, the needle holes are connected with each other to form a continuous slit, thereby making it impossible to transplant hair segments onto a base.\nMore importantly, when the needle is to penetrate the base, it pushes the base, so that the needle would move and wobble.\nIf the needle should wobble, the hair segment cannot surely be hooked by the needle.\nAfter repeated trial and error in development of automated wig manufacturing apparatus which is the first in the world, the present inventor has reached a conclusion that a keyword is an issue of the needle wobbling.\nThe present invention has been made in view of the above-described background, with the object to surely hook the hair segment with the needle in automated wig manufacturing. Another object is to control the transplanting pitch of the hair segment with great accuracy. Still another object is to reduce a percentage of production of defective articles when automatically manufacturing wigs."}
{"text": "This invention relates to cotton bale covers and, more particularly, to a circular knit cotton cover designed to be used on high density cotton bales.\nA wide variety of covers for cotton bales have been used in the past. Those made from jute or burlap are loosely woven permitting contamination of the bale to occur and their low strength causes rips and tears of the cover to develop during handling and storage. Attempts have been made to construct covers of other materials to avoid contamination and tearing. In U.S. Pat. No. 3,647,061 the bale cover disclosed is a laminate of a paper outer layer and a nonwoven cotton inner layer bonded through a polypropylene net with a low melt acrylate. Such construction is expensive and its low air permeability tends to cause mildew. A similar nonwoven cotton laminate cover is shown in U.S. Pat. No. 3,647,139 wherein two nonwoven cotton layers are bonded in a like manner. These covers are also very expensive in comparison to jute and have the further disadvantage of being non-biodegradable. An attempt was then made to produce a knitted cotton cover in U.S. Pat. No. 4,071,138. This bale cover was made from a warp knit fabric having tricot stitch construction. However the bag must be shaped by spiral sewing which adds to its weight, increases cost and reduces its stretch due to the tricot construction and the spiral seams. Synthetic fibers were next used in cotton bale covers. U.S. Pat. No. 4,557,985 teaches the use of a woven polyolefin fabric with resin coating stripes to help reduce fraying as a cotton bale cover. This cover bag has good strength but is non-biodegradable, heavy and expensive to make since it has to be spiral-tubed. An in U.S. Pat. No. 5,104,703 the cotton bale fabric is a single layer nonwoven fiber blend of cotton and polypropylene. Here also the product is not biodegradable and is prepared by spiral tube sewing with reduced stretch in comparison to cotton covers.\nThe cotton bale cover bag of this invention is completely biodegradable, can easily be recycled due to its jersey stitch construction, has improved stretch and shape conformity characteristics, is light weight, has excellent strength and tear resistance since it has only one seam across the end or bottom portion as it is not spiral sewn and is economically attractive."}
{"text": "Machines in the prior art such as copiers have cabinetry and housings which are generally flat, such as flat access panels and flat top, front, and side surfaces. Flat panels and surfaces are economical to manufacture, and may be formed with the components of the machine to constitute a relatively compact configuration.\nHowever, flat panels give no immediate indication of the functions of the trays and panels. Paper tray access panels are typically indistinguishable from regions presenting information for visual output to the user. That is, input regions and output regions are not differentiated by shape, as all of such regions are generally flat. Accordingly, labels with indicia thereon are typically required to be provided and affixed to the appropriate panels to indicate the function of each panel, which thus increases the assembly time and cost of the machine. In addition, users and service personnel must learn the use of the panels by instructions and/or by experience. Such learning reduces the efficiency of the users and service personnel in having to learn and memorize the functions of each access panel as needed.\nAccordingly, a need exists for cabinetry and housings of machines which provides quick visual indications of the functions of access panels."}
{"text": "1. Field of the Invention\nThe present invention relates to a quick-action flap according to the preamble of claim 1.\n2. Discussion of Background\nA quick-action fitting together with an intercept valve are placed in steam turbines between the respective reheater and the intermediate pressure turbine. In dependence on the steam data and the type of mounting, the quick-action fitting can be a valve or a flap, preferably a swing flap. The opening of this quick-action fitting occurs hydraulically, closing is initiated via a spring. A swing flap possesses the advantage over a valve that it does not have to deflect the steam and therefore, naturally causes lower flow losses. Furthermore, a swing flap facilitates an advantage arrangement on the intermediate pressure turbine. On the other hand, it is to be noted that a swing flap is subjected to this flow to a greater degree than a valve and thus experiences greater mechanical loading, i.e., in respect of loading, the swing flap structure is highly stressed. If the flow processes in a swing flap are considered, it is to be noted that, during the closing process, the steam pressure on the outflow side of the flap falls and thus the pressure difference at the disk of the flap rises. The entire swing flap is speeded up by the force growing in the flow direction to such an extent that a high, and for the construction of the swing flap, dangerous impact speed of the flap disk against the seating ring in the flap casing occurs. Although the conclusion could be drawn that the construction should be strengthened, with increased mass the kinetic energy would also rise. The impact speed could be reduced by means of a hydraulic brake; however, any deceleration causes braking forces (mass times deceleration) which would greatly load the entire mechanism of the swing flap. It is also to be borne in mind that, even if the mechanism of the flap withstands the loading caused by the braking force, a hydraulic brake designed for a low closing speed would have the great disadvantage for such a flap that the closing speed would be too low in respect of the task and correspondingly the closing time would be too long, not to mention the fact that, in such a constellation, the steam pressure behind the flap would fall even further, which would cause the pressure difference and thus the force on the flap to rise even more.\nIn addition, the smaller the chambers through which there is flow, between the quick-action flap and the entry to the blading, the quicker the pressure behind the flap disk must drop due to falling accumulation effect, which in turn must have a negative effect on the impact speed of the flap against the seating ring. If this were to be remedied, one would have to provide long lines in each case behind the flap, which alone could provide large accumulation capacities. However, it is part of the definition of the object to place the fittings directly in front of the turbine;\nin order to minimize the accumulation capacity, in order to avoid damage to the turbine due to excessive speed of rotation; PA1 because of more favorable use of space."}
{"text": "Plastic and reconstructive surgery relies greatly on the use of free skin grafts, expansion of contiguous tissue or pediculed flaps of various shapes and sizes. In the repair of skin, the amount of skin available is often not sufficient for grafting and reconstruction. The requirement for additional skin or tissue has led to the development of various types of tissue expander devices. These devices are generally designed to cause gradual expansion of the skin and its subcutaneous tissue by applying progressively increasing pressure over a predetermined time period. An example of naturally occurring tissue expansion is seen in the extension of the skin over the pregnant female abdomen. A similar condition occurs when obese individuals undergo weight loss resulting in the development of folds of excess tissue.\nTypically, the expander devices now commercially available and in use employ some form of elastic envelope or balloon which is surgically implanted under the skin and subcutaneous tissue in the area of the desired expansion. Ordinarily, a reservoir and connecting tube are implanted under the subcutaneous layers in communication with the envelope, and the envelope is slowly inflated by injecting a sterile saline solution into the reservoir using a syringe and hypodermic needle. The solution passes from the reservoir into the envelope through the connecting tube so as to inflate the envelope and cause it to exert pressure on the surrounding tissue which in turn expands to accommodate the increase in volume of the envelope. Periodic sterile injections of additional saline are typically performed at intervals of several days apart, this process being repeated until the desired degree of expansion has been achieved. When expansion has been completed, the envelope is surgically removed and the expanded tissue is then available for the desired plastic or reconstructive surgery. Representative tissue expander devices are set forth and described in a number of patents including U.S. Pat. No. 4,217,889 to C. Radovan; U.S. Pat. No. 4,666,447 to G. M. Smith et al; U.S. Pat. No. 4,615,704 to E. E. Frisch; U.S. Pat. No. 4,651,717 to E. R. Jakubczak; U.S. Pat. No. 4,662,357 to D. L. Pierce et al and U.S. Pat. No. 4,685,447 to A. A. Iversen et al.\nNotwithstanding the development and availability of numerous types of tissue expander devices, a problem which is especially prevalent in connection with the use of these devices is that of closely controlling expansion of the tissue to avoid undue pain as well as tissue necrosis. For example, the aforementioned patent to Radovan et al describes the problems associated with too rapid expansion of tissue and which may lead to lack of blood flow in the surrounding expanded tissue which if not corrected results in tissue necrosis and loss of tissue viability. Similarly, wound dehissence may occur requiring a further healing period before expansion can be resumed. Presently, in the course of each filling or inflation procedure, the envelope is inflated until the tissue becomes \"tight\" or \"firm\" and this degree of tightness is based on the subjective opinion of the clinician; or, the clinician may gauge the degree of inflation by monitoring the pain experienced by the patient. However, this is not a reliable method as patients have widely variable thresholds of pain.\nPressure measuring devices have been utilized in the past for monitoring fluid pressure in the body. For example, U.S. Pat. No. 4,571,749 and U.S. Pat. No. 4,731,083 to R. E. Fischell disclose pressure measuring devices for a urinary sphincter. Similarly, U.S. Pat. No. 3,744,063 to D. M. McWhorter et al is directed to a method and apparatus for pressure monitoring an artificial sphincter. U.S. Pat. No. 3,863,584 to A. S. Borsanyi is concerned more with measuring the fluid pressure of plasma, blood or other fluids in the body while isolating the gauge from direct contact with the fluid. To the best of my knowledge, no one has recognized the ability to monitor tissue expansion by providing a reproducible method of measuring the pressure exerted on the tissue during or after the filling procedure. By closely monitoring the pressure and providing for a precise means of measuring actual pressure, it is possible to achieve an optimum inflation rate or degree of inflation which will overcome and avoid the previously noted problems associated with tissue expander devices. In this connection, it is hiqhly desirable that the pressure monitor be capable of directly sensing the fluid or liquid pressure condition of the tissue expander envelope but at the same time be isolated from the liquid such that the sterility of the fluid path is not compromised. That is to say that the liquid contents of the expander must be protected from any bacteria or particulate matter that may be carried by the liquid to or from the monitor itself. In this way, even though the tissue expander fill kit and tissue expander are disposable, the pressure gauge or monitor may be reusable."}
{"text": "1. Field of the Invention\nThe invention relates to a high density optical disk on which recording and reproduction of information are conducted with respect to a land and a groove of a track, and also to an information recording/reproduction apparatus which conducts recording and reproduction of information on the optical disk.\n2. Description of the Prior Art\nAn optical disk is a high density memory which has a large capacity, on which non-contact recording and reproduction can be conducted, and which can be replaced with another one. Capacities of typical optical disks are as follows: When an optical head of a lens NA of 0.5 is used, a 130-mm disk has a capacity of 300 to 500 MB per face, and a 90-mm disk has a capacity of about 128 to 250 MB per face. For a multimedia purpose, a high density recording/reproduction technique which uses a short-wavelength laser of 680 nm to obtain a capacity that is about 2 to 4 times the above-mentioned capacities has been studied.\nFIG. 10 shows plan views (upper ones) and section views (lower ones) of continuous servo track formats in the prior art.\nFIG. 10(a) shows a continuous servo track which is employed in a 130-mm or 90-mm optical disk of the prior art and which is of a land record track format. In the land record track format, a track consists of a groove 2 which is formed on a transparent substrate 1 and which has a depth of .lambda./(8.multidot.n) (.lambda. is the wavelength of a laser, and n is the refractive index of the substrate 1. The same applies to the followings.), and pits 4 constituting a sector identification (ID) signal, and record marks 5 are recorded onto a land 3 sandwiched by the tracks. The pits 4 for the ID signal are convex-concave pits having a phase depth of .lambda./(4.multidot.n).\nThe track pitch is selected so as to be about .lambda./NA which is obtained from the laser wavelength (.lambda.) and a lens aperture (NA). In the prior art disk, since the lands must remain to exist between the grooves 2 and the ID signal pits 4, it is difficult to reduce the track pitch to 1.3 .mu.m or less in the view point of a land forming process.\nFIG. 10(b) shows an example of a groove record track format in which tracks are formed by simple grooves 6 having a phase depth of .lambda./(8Nn) and lands 7, and pits 8 constituting an ID signal and record marks 9 wherein data signals are recorded are recorded into the grooves 6. Since such a groove record track has a track structure consisting of the simple grooves 6, a disk of a track pitch which is 1 .mu.m or less can easily be produced.\nFIG. 10(c) shows an example of a land/groove record track format in which groove record tracks are formed by setting the width of grooves 11 having a depth of about .lambda./(8.multidot.n) to be one-half of the track pitch, and signals 12 are recorded also onto lands 10. In principle, this land/groove recording can achieve an area recording density which is twice that of the land recording of FIG. 10(a).\nGenerally, when a track pitch is to be reduced, there arise problems in a cross-talk of a signal recorded in an adjacent track, in a cross-erase in which also signals into two adjacent tracks are erased by conducting a data recording, and in the stability of a tracking servo.\nThe stability of a tracking servo will be discussed. In the land record track 10 for the land/groove recording of FIG. 10(c), the track pitch is one-half of .lambda./NA. In the case where .lambda. is 830 nm and NA is 0.5, even when the pitch of tracks for recording signals is 0.8 .mu.m, the tracking servo can be conducted with respect to the track pitch of 1.6 .mu.m consisting of each groove and land. Consequently, the tracking can be done stably by either of the 3-beam method and the push-pull method which are conventionally employed.\nEven in the land/groove recording in which recording is conducted as described above, when the track pitch is further reduced in order to increase the recording density, there arises a problem in that a cross-talk between the track of the groove 11 and that of the land 10 occurs.\nSpecifically, in the case where an optical head in which the laser wavelength .lambda. is 830 nm and the lens aperture NA is 0.5 is used, a cross-talk of -30 to -35 dB occurs when the track pitch is 1.6 .mu.m, and that of -15 to -20 dB occurs when the track pitch is 0.8 .mu.m, thereby producing a problem in that an ID signal and a data signal cannot be reproduced normally.\nParticularly, there arise the following problems: In a track search process, the confirmation of a target track is made difficult by reproduction errors such as that an ID signal is caused to be erroneously reproduced by an effect of a cross-talk on the ID signal, and that an ID signal of an adjacent track is erroneously reproduced. In a recording process conducted on an unrecorded sector, even when an ID signal leaks from an adjacent track, the reproduction is apparently conducted in a normal manner, thereby causing data to be recorded into a wrong sector.\nIn the groove/land recording conducted as described above, however, the track density is doubled, and hence there arises a problem in that, in a recording or reproducing process, a cross-talk, or a cross-erase produces errors which exceed a criterion, so that defective sectors are generated in an increased number."}
{"text": "Clearing of obstructions from drain pipes typically involves hiring a plumber or technician to clean out the obstruction using a drain auger, especially when the obstruction is quite substantial and consumer grade hand-powered drain augers are insufficient for the task, as the typical homeowner is unlikely to own or operate a heavier grade motor-driven unit. The need for professional services not only increases the cost involved in a drain cleanout operation, but can also lead to significant wait times for the consumer depending on the availability of services providers in the area. Customers may also be unlikely to schedule regular maintenance in view of the cost and/or inconvenience of seeking a professional service provider, thereby increasing the likelihood of more significant obstructions building up over time and leading to potential drain line backup, and associated property damage.\nU.S. Pat. No. 4,376,321 discloses an automatic drain cleanout device in the form of a motorized, timer activated drain snake with a suitable fitting for attachment to a plumbing adapter between the drain trap of a sink and a vertical drain pipe. While this provides for improved plumbing maintenance via regular ongoing cleanout operations, this device lacks any means for retracting the drain snake, which therefore remains within the plumbing path at all times, reducing the available flow area within the pipes.\nU.S. Pat. No. 4,388,741 teaches an automatic sewer cleaning system that similarly prevents clogs through regular timer-controlled cleanout operations, but again uses a cleanout cable that remains within the plumbing at all times.\nAccordingly, there remains room for improved solutions in the field of automated drain cleanouts, in response to which Applicant has developed a unique drain cleanout apparatus."}
{"text": "This invention relates generally to a heat transfer labeling apparatus and a method of applying a heat transferable label to a hollow article. More particularly, this invention relates to such an apparatus and method wherein the articles are decorated in a continuous operation.\nNumerous decorating techniques are known in the art, some of which include the application of a label onto a hollow article to be decorated. One of the techniques which is desirable in this type of decorating is the usage of a heat transferable label which includes a decorative predetermined design thereon and may thus be transferred onto the article or container being decorated.\nThe heat transfer process permits for multicolored designs to be applied to a container in a single operation. The heat transfer process involves the use of a release-coated carrier upon which the design to be transferred is printed. The design is transferred from the web-like carrier to the container generally by using a combination of heat and pressure. The principal advantage of the heat transfer technique is that multicolored designs of an infinite variety may be applied to a container.\nBecause of the heat requirements associated with the release and application of the label from the web onto the container, it has been generally accepted practice to maintain the container in a stationary position, albeit rotatable in the instances of circular containers, during the decorating step. This has resulted in numerous prior art types of apparatus which employ intermittently moving mechanisms which include one to engage and deposit a container at a decorating station. Yet another mechanism engages the container at a decorating station. This latter mechanism must permit relative movement between the container and die to facilitate application of the label onto the container to be decorated. Once the container has been decorated it is removed from the decorating station by yet another mechanism and conveyed to another destination. Each of these functions has required numerous types of moving parts and mechanisms to impart the desired motion and transfer of the container to and from a decorating station.\nBecause of the intermittent movement associated with such systems, the speed of decoration has been limited. The various movements have curtailed operating speeds and placed heat transfer labeling systems in a limited and low rate of production category. Because the means disclosed herein have eliminated these various undesirable mechanisms, a system has been devised which overcomes the disadvantages of the prior art and provides for a system wherein articles are decorated in a continuous manner."}
{"text": "The present invention relates to plastic compositions containing a polysulfone. More specifically, the present invention covers plastic compositions containing a three component system which consists of a polysulfone, certain bis-phenoxy compounds (hereinafter defined) as flame retardants for said plastic compositions, and an enhancing agent (hereinafter defined) for said flame retardants.\nPolysulfones and utility thereof are known in the art as exemplified by Encyclopedia of Polymer Science and Technology Vol. 9, pages 460-485, John Wiley & Sons, Inc., New York, 1968, and Modern Plastics Encyclopedia 1972-1973, Vol. 49, No. 10A, October, 1972, pages 98, 102, 162 and 222 and which publications are in toto incorporated herein by reference.\nThe need for flame retarding polysulfones has also been recognized in the art as exemplified by Modern Plastics Encyclopedia, ibid, pages 222 and 456-458.\nThe prior art has specifically recognized the problems of finding suitable flame retardants for polysulfones in view of the fact that polymer systems differ substantially in both flammability characteristics and physical properties and there is no predictability whatsoever from one system to another. Thus in the Norris et al paper entitled \"Toxicological and Environmental Factors Involved in the Selection of Decabromodiphenyl Oxide as a Fire Retardant Chemical\", Applied Polymer Symposium No. 22, 195-219 (1973), the authors state: \"A growing recognition of the huge annual toll taken by fire is resulting in more stringent flammability requirements for synthetic polymers in a variety of applications. Because of economic constraints and the need to produce flame resistant polymers without total replacement of existing manufacturing processes, increased flame resistance is generally achieved by incorporation of a fire retardant chemical in the finished product. This chemical is usually based on bromine, chlorine, phosphorus, or nitrogen and may either be chemically reacted or physically blended into the product. Since polymer systems differ markedly in both flammability characteristics and physical properties, selection of a suitable flame retardant depends on a variety of factors that severely limits the number of acceptable materials.\n\"A general class of synthetic polymers that require flame retardancy because of their use in electrical and high temperature applications, but pose severe problems in selecting a suitable flame retardant are the high performance thermoplastic resins such as thermoplastic polyesters, polyphenylene oxides, and acrylonitrile-butadiene-styrene (ABS) terpolymers. Some of the most important criteria for an acceptable flame retardant in these applications are:\n1. It must be as effective as possible to minimize both cost and effect on polymer properties. Use levels may range up to 15% by weight.\n2. It must have sufficient stability to withstand conditions encountered during polymer processing and use. Processing conditions (blending, extrusion, and molding) often involve temperatures exceeding 300.degree. C. The flame retardant must tolerate these conditions without degradation or volatilazation. Also, attention must be given to hydrolytic stability and oxidative degradation, particularly under extended service at high temperatures.\n3. It must be compatible with the base polymer and exert minimal adverse effect on those properties that give the polymer its value. Some of these critical properties are tensile strength, impact strength, heat deflection temperature, shear strength, and flexural modulus.\n4. Finally, the flame retardant must not interfere with attainment of desired product esthetics and form.\n\"Because of the stringest thermal stability requirements, only a very few compounds have been identified which can meet the necessary performance and economic criteria.\"\nThe resultant disadvantages in the utilization of various prior art materials as flame retardants, in general, for plastic compositions include, without limitation, factors such as thermal migration, heat instability, light instability, non-biodegradable, toxicity, discoloration, the large amounts employed in order to be effective, and the unpredictable end results obtained when using the same material in different plastics (note, for example, in Modern Plastics Encyclopedia, ibid, page 650, wherein octabromobiphenyl is suitable for use in polyolefins as a flame retardant therefor, but is not shown for use (or functionally equivalent) as such for the other 27 compositions listed such as ABS; polycarbonates, polystyrene, acrylics and polyurethanes). Thus, it can be seen that the field of flame retardancy is highly sophisticated and consequently requires substantial research effort to achieve a particular desired end result.\nIn conjunction with the rendition of a polysulfone flame retardant the aforegoing discussion is particularly applicable. Furthermore, it is desirable that, in addition to the retention of good physical characteristics, three characteristics (one of which is critical) of the product (polysulfone) be within certain defined limitations in order to provide a functional flame retarded product. These three characteristics are (1) light stability (as measured, for example, by .DELTA. E color values, hereinafter defined), (2) flame retardancy (as measured, for example, by UL-94, hereinafter defined, and which is a critical characteristic and limitation), and (3) thermal stability (as measured, for example, by certain ASTM tests for decomposition and migration, hereinafter defined). This characteristic of flame retardancy is not an arbitrary item but is a criteria which is strictly adhered to in the polysulfone art and hereinafter explained in more detail, with particular reference to the examples contained herein.\nIn conjunction with the foregoing discussion, the prior art in general suggests the use of halogen-containing materials as \"potential\" or \"possible\" flame retardants for plastic materials. However, the prior art also recognizes that any material must be adjudged on a case by case basis because of the unpredictable results of the end product when any additive is incorporated therein. For example, with reference to the use of a halogenated fire retardant in U.S. Pat. No. 3,658,634 attention is directed to the fact that the patentee specifically points out the disadvantages in the use of a halogen-containing fire retardant. In Column 1, lines 14-17, the patentee states: \"Therefore, if it is possible to impart fire retardancy to the thermoplastic polymers without deteriorating the useful properties of the thermoplastic polymers, they can be widely used in the field of inertia, construction and electric industries.\" In Column 1, lines 26-32 the patentee states: \"--the compounds containing chlorine or bromine atoms to be used as fire-retardant agents, are generally sublimated and therefore, the fire retardant agents are sublimated and lost in the process for producing fire-retardant polymers or in after-finishing processes; accordingly, deteriorations of fire-retardancy or difficulties in use tend to occur more often than not.\"\nIn Column 1, lines 39-44 the patentee states: \"--the compounds containing chlorine or bromine atoms to be used as fire-retartant agents are unstable in most cases when exposed to ultraviolet rays.\" In Column 1, lines 59-64 the patentee states: \"However, as a matter of fact, only very few fire-retardant polymers can be used in actual practice although they are said to have fire-retardant effects, because there are restrictions such as the conditions employed in product attributable to the properties of the fire-retardant agent, or to the properties of the polymers into which they are to be incorporated.\"\nIt can be seen, then, from the foregoing discussion and quoted subject matter that the field of flame retardancy is highly sophisticated, unpredictable and requires substantial research to produce an end product (plastic composition) which meets the necessary criteria for utilitarian purposes, particularly under the present day government standards. Thus, there is always a demand for a material which will function as a flame retardant in a polysulfone and concurrectly will not, be incorporation therein, adversely effect the chemical and/or physical and/or mechanical properties of the resultant polysulfone containing plastic composition (herein also referred to as \"polysulfone plastic composition\") and also have utility.\nThe prior art problem of providing a flame retarded polysulfone composition having desired chemical, physical and mechanical properties, in addition to functional utility, has now been substantially solved by the present invention and the above-described disadvantages substantially overcome.\nAccordingly, one of the main objects of the present invention is to provide polysulfone plastic compositions which are flame retarded.\nAnother object of the present invention is to provide a unique three component system for polysulfone plastic compositions which will not substantially adversely affect the chemical and/or physical and/or mechanical properties of said compositions.\nA salient object of the present invention is to provide a polysulfone plastic composition which has a certain defined flame retardancy property.\nA further object of the present invention is to provide a flame retardant and an enhancing agent which are economical and easy to incorporate into polysulfones without being degraded or decomposed as a result of normal blending or processing operations."}
{"text": "The rapid proliferation of portable electronic devices in the international marketplace has led to a corresponding increase in the demand for advanced secondary batteries. The miniaturization of such devices as, for example, cellular phones, laptop computers, etc., has naturally fueled the desire for rechargeable batteries having high specific energies (light weight). Mounting concerns regarding the environmental impact of throwaway technologies, has caused a discernible shift away from primary batteries and towards rechargeable systems.\nIn addition, heightened awareness concerning toxic waste has motivated efforts to replace toxic cadmium electrodes in nickel/cadmium batteries with the more benign hydrogen storage electrodes in nickel/metal hydride cells. For the above reasons, there is tremendous market potential for environmentally benign secondary battery technologies.\nSeveral approaches have been pursued in an effort to develop improved secondary battery technologies, including the recent introduction of the nickel/metal hydride cell and the commercialization of lithium-ion technologies. Among the factors leading to the successful development of high energy density batteries, is the fundamental need for high voltage and/or low equivalent weight electrode materials. Electrode materials must also fulfill the basic electrochemical requirements of sufficient electronic and ionic conductivity, high reversibility of the oxidation/reduction reaction, as well as excellent thermal and chemical stability within the temperature range for a particular application. Importantly, the electrode materials must be reasonably inexpensive, widely available, non-toxic, and easy to process.\nIt has been found previously that organodisulfide compounds can be used as high energy density electrodes for rechargeable battery systems. In the work by Liu et al. [Liu et al., J. Electrochem. Soc., 138: 1891 (1991); Liu et al., J. Electrochem. Soc., 138: 1896 (1991); Visco et al., Mol. Cryst. Lig. Cryst. 190: 185 (1990); and Visco et al., U.S. Pat. No. 5,162,175 (issued Nov. 10, 1992)], a novel class of polymeric organodisulfides were described having exceptionally low equivalent weights and consequently very high gravimetric and good volumetric capacities. The electrochemical reaction involves the oxidation of a thiolate anion to a sulfur radical which rapidly dimerizes to form a disulfide linkage as shown below: EQU 2[.sup.- SRS.sup.- -e.sup.-.fwdarw..sup.- SRS ] EQU 2.sup.- SRS .fwdarw..sup.- SRS-SRS.sup.-.\nThose reactions lead ultimately to the formation of a polymeric organodisulfide.\nIn the preferred embodiment, the polyorgano-disulfides (PDS) were used as positive electrodes in thin-film polymer electrolyte cells having alkali metal negative electrodes. [See, Visco et al., U.S. Pat. No. 5,162,175.] Since organodisulfides are invariably electronic insulators, carbon black was included in the electrode formulation. The PDS were formulated as an intimate mixture of the disulfide, solid polymer electrolyte (SPE), and carbon black. Although the intermediate temperature performance (60.degree. C. to 120.degree. C.) of the alkali metal solid polymer batteries having PDS cathodes is excellent, utilization of positive electrode capacity can be critically dependent on the intimate connection between the carbon black particles, PDS particles, and current collectors. That interconnectivity appears to be more important as the temperature of operation is reduced to ambient and sub-ambient temperatures.\nThe instant invention provides a means to overcome the temperature range limitations of the Liu et al. PDS electrodes. The need for a high degree of interconnectivity between and among the carbon black, PDS matrix, and current collectors of the Liu et al. PDS electrodes is eliminated or minimized according to the instant invention by the introduction of metallic conductivity into the polyorganodisulfide polymer chain.\nMetallic conductivity and low temperature superconductivity has been reported for inorganic polymers such as (SN).sub.x, and synmetals such as tetracyano-p-quinodimethane-tetrathiafulvalene (TCNQ-TTF) since the early 1970s. In the mid-1970s considerable attention was directed towards the study of organic polymers following the discovery that polyacetylene could be prepared as a film having metallic luster.\nIn the late 1970s it was found that the conductivity of polyacetylene could be increased 13 orders of magnitude by doping it with various donor or acceptor species, reaching conductivities of 10.sup.3 S.multidot.cm.sup.- 1. It was later proposed by Heeger and MacDiarmid [MacDiarmid et al., J. Phys. Collog. 44: C3-543 (1983)] to use such materials as electrodes in secondary batteries. Unfortunately, with few exceptions, conducting polymers such as doped polyacetylene tend to be air sensitive, difficult to process, and have very low volumetric capacities. For those reasons, most groups have abandoned conducting polymers for battery applications.\nRecently, highly conducting metallo-organic and/or coordination polymers have been reported which exhibit exceptionally good electrical properties and environmental stability, as well as multiple oxidation states. For example, tetrathiolate ligands can be reacted with transition metal salts (see reaction set forth below) to yield black amorphous coordination polymers with exceptionally high electrical conductivity and environmental stability [Dahm et al., Synthetic Metals, 55-77: 884-889 (1993)]. Such materials have achieved electrical conductivities of 10.sup.2 S.multidot.cm.sup.-1. ##STR1##\nThe instant invention concerns the use of metallo-organic charge-transfer materials as positive electrodes in secondary batteries. Some metal oxides and chalcogenides presently used as battery electrodes have appreciable electrical conductivities; however, many such oxide/chalcogenide materials need to be formulated with carbon black to increase the conductivity of the composite electrode to acceptable levels. In contrast, the inherently high electrical conductivity of charge-transfer materials allows them to be used without significant dilution.\nThe metallo-organic charge-transfer materials of this invention are essentially redox electrodes and in principle are reversible to an enormous variety of counter electrodes. Oxide/chalcogenide cathode materials presently under investigation by other groups are not reversible to alkali metal ions other than lithium, and thus, lack the flexibility of the electrodes of the instant invention.\nThe reversibility of the electrodes of the instant invention can furnish large economic advantages in the commercialization of secondary batteries based on charge-transfer materials, in that expensive lithium anodes can be replaced by inexpensive sodium electrodes. Also, the ability to alter the thermodynamic redox potential of the charge-transfer electrode by a suitable choice of the metal ion and/or of the organosulfur chelating ion provides tremendous flexibility in tailoring-the secondary battery characteristics to the specific application at hand.\nTeo [U.S. Pat. No. 4,181,779 (issued on Jan. 1, 1980)] describes the use of halogen reactive materials such as organo-metallic polymers limited to the nominal stoichiometry [M(TTL)].sub.x in which M is a transition metal containing complex having at least one transition metal ion selected from Group VIII of the periodic table; wherein TTL has the nominal atom composition C.sub.10 H.sub.4 X.sub.4 and substituted compositions thereof in which X is selected from the group consisting of sulfur (S), selenium (Se), tellurium (Te), and mixtures thereof; and x is equal to or greater than 1. The TTL compounds described by Teo have high formula weights, for example, as shown below, leading to unattractive energy densities (watthours/kilogram) for electrodes formulated with such materials. ##STR2##\nIn contrast to the Teo electrodes, the electrodes of the instant invention have low equivalent weights and consequently high energy densities.\nA key feature of the electrodes of the instant invention is the use of thiolate ligands to chelate metal ions of the electrodes and thereby form coordination polymers having low equivalent weights and enhanced electrical conductivity. The chelation also solves a prior art dissolution problem of the PDS electrodes. For example, thiolate anions from the PDS electrodes most notably in a liquid or gel format, diffuse and migrate to the negative electrode, which may result in deterioration of battery performance. In accordance with the instant invention, the metal-organosulfur polymers are anchored by chelation and locked into position.\nThe new metal-organosulfur positive electrodes of this invention, thus, overcome as indicated above many of the problems found in prior art battery systems. The invention provides for novel secondary batteries that have high electrical conductivity, high energy densities and polymeric positive electrodes that are stably situated. Further, the secondary cells of this invention represent improvements over prior art batteries in having more cycle life and in being able to perform at a lower temperature without deterioration of performance."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of downhole sampling and in particular to the maintenance of hydrocarbon samples in a single-phase state after capture in a sample chamber.\n2. Summary of the Related Art\nEarth formation fluids in a hydrocarbon producing well typically comprise a mixture of oil, gas, and water. The pressure, temperature and volume of formation fluids control the phase relation of these constituents. In a subsurface formation, high well fluid pressures often entrain gas within the oil above the bubble point pressure. When the pressure is reduced, the entrained or dissolved gaseous compounds separate from the liquid phase sample. The accurate measure of pressure, temperature, and formation fluid composition from a particular well affects the commercial interest in producing fluids available from the well. The data also provides information regarding procedures for maximizing the completion and production of the respective hydrocarbon reservoir.\nCertain techniques analyze the well fluids downhole in the well bore. U.S. Pat. No. 6,467,544 to Brown, et al. describes a sample chamber having a slidably disposed piston to define a sample cavity on one side of the piston and a buffer cavity on the other side of the piston. U.S. Pat. No. 5,361,839 to Griffith et al. (1993) disclosed a transducer for generating an output representative of fluid sample characteristics downhole in a wellbore. U.S. Pat. No. 5,329,811 to Schultz et al. (1994) disclosed an apparatus and method for assessing pressure and volume data for a downhole well fluid sample.\nOther techniques capture a well fluid sample for retrieval to the surface. U.S. Pat. No. 4,5 83,595 to Czenichow et al. (1986) disclosed a piston actuated mechanism for capturing a well fluid sample. U.S. Pat. No. 4,721,157 to Berzin (1988) disclosed a shifting valve sleeve for capturing a well fluid sample in a chamber. U.S. Pat. No. 4,766,955 to Petermann (1988) disclosed a piston engaged with a control valve for capturing a well fluid sample, and U.S. Pat. No. 4,903,765 to Zunkel (1990) disclosed a time delayed well fluid sampler. U.S. Pat. No. 5,009,100 to Gruber et al. (1991) disclosed a wireline sampler for collecting a well fluid sample from a selected wellbore depth, U.S. Pat. No. 5,240,072 to Schultz et al. (1993) disclosed a multiple sample annulus pressure responsive sampler for permitting well fluid sample collection at different time and depth intervals, and U.S. Pat. No. 5,322,120 to Be et al. (1994) disclosed an electrically actuated hydraulic system for collecting well fluid samples deep in a wellbore.\nTemperatures downhole in a deep wellbore often exceed 300 degrees F. When a hot formation fluid sample is retrieved to the surface at 70 degrees F., the resulting drop in temperature causes the formation fluid sample to contract. If the volume of the sample is unchanged, such contraction substantially reduces the sample pressure. A pressure drop changes in the situ formation fluid parameters, and can permit phase separation between liquids and gases entrained within the formation fluid sample. Phase separation significantly changes the formation fluid characteristics, and reduces the ability to evaluate the actual properties of the formation fluid.\nTo overcome this limitation, various techniques have been developed to maintain pressure of the formation fluid sample. U.S. Pat. No. 5,337,822 to Massie et al. (1994) pressurized a formation fluid sample with a hydraulically driven piston powered by a high-pressure gas. Similarly, U.S. Pat. No. 5,662,166 to Shammai (1997) used a pressurized gas to charge the formation fluid sample. U.S. Pat. No. 5,303,775 (1994) and U.S. Pat. No. 5,377,755 (1995) to Michaels et al. disclosed a bi-directional, positive displacement pump for increasing the formation fluid sample pressure above the bubble point so that subsequent cooling did not reduce the fluid pressure below the bubble point.\nExisting techniques for maintaining the sample formation pressure are limited by many factors. Pretension or compression springs are not suitable because the required compression forces require extremely large springs. Shear mechanisms are inflexible and do not easily permit multiple sample gathering at different locations within the well bore. Gas charges can lead to explosive decompression of seals and sample contamination. Gas pressurization systems require complicated systems including tanks, valves and regulators which are expensive, occupy space in the narrow confines of a well bore, and require maintenance and repair. Electrical or hydraulic pumps require surface control and have similar limitations.\nAccordingly, there is a need for an improved system capable of compensating for hydrostatic well bore pressure loss so that a formation fluid sample can be retrieved to the well surface at substantially the original formation pressure, that is, in a single phase state. The system should be reliable and should be capable of collecting the samples from the different locations within a well bore.\nUnlike an ordinary sample tank, however, a single-phase tank has a floating piston inside of it. Sample fluid or crude is pumped into the sample tank against the top side of the piston. Downhole, as crude oil is pumped into the tank, the pumped crude pushes against the top side of the floating piston inside of the sample tank and further compresses the gas cushion underneath the sample tank piston. Crude oil is pumped into the sample tank against the cushioned piston until its pressure is several thousand pounds per square inch above formation pressure. The gas cushion is initially created at the surface where the tank is charged before going into the well bore. The purpose of charging the down hole sample tank is to maintain the down hole sample of crude oil in a single phase condition after it has been brought to the surface and cools. Gas is pumped underneath the sample tank piston to charge the sample tank cylinder.\nTo charge the single-phase sample tank cylinder a non-reactive gas (e.g., nitrogen) is connected to the sample tank through a pressure regulator. The tank is filled until the pressure underneath the sample tank piston reaches the set pressure of the regulator. The tank inlet valve is then closed thereby trapping as many moles of gas as can possibly fit into the tank volume underneath the piston at that pressure. This gas cushion is important when collecting down samples of crude oil at elevated temperatures of 100-200 C and pressures of 10-20 kpsi. As these tanks are brought back to the surface, the tank and the sample inside of the tank, once removed from the high temperature down hole in the well bore, cools to the ambient surface temperature so the crude oil within the sample tank shrinks or reduces its volume and pressure associated therewith is likewise reduced. This temperature-induced shrinkage can be as much as 30% of the initial crude oil volume. At this reduction in pressure, below the bubble point for the crude, it is expected that natural gas bubbles will nucleate or asphaltenes precipitate and come out of the crude oil and fill the void left by shrinking liquid. Nucleation of gas bubbles or precipitation of solids changes the single-phase liquid crude to a two-phase state consisting of liquid and gas or liquid and solids. Two-phase samples are undesirable, because once the crude oil sample has separated into two phases, it can be difficult or impossible and take a long time (weeks), if ever, to return the sample to its initial single-phase liquid state even after reheating and/or shaking the sample to induce returning it to a single-phase state.\nDue to the uncertainty of the restoration process, any pressure-volume-temperature (PVT) lab analyses that are performed on the restored sing-phase crude oil are often suspect. When using ordinary sample tanks, one tries to minimize this problem of cooling and separating into two-phase by pressurizing the sample down hole to a pressure that is far (4500 or more psi) above the downhole formation pressure. The extra pressurization is an attempt to squeeze enough extra crude oil into the fixed volume of the tank that upon cooling to surface temperatures the crude oil is still under enough pressure to maintain a single-phase state and maintains at least at the pressure that it had downhole.\nThe gas cushion of the single-phase tanks, thus, makes it easier to maintain a sample in a single phase state because, as the crude oil sample shrinks, the gas cushion expands to keep pressure on the crude. However, if the crude oil shrinks too much, the gas cushion (which expands by as much as the crude shrinks) may expand to the point that the pressure applied by the gas cushion to the crude falls below formation pressure and allows asphaltenes in the crude oil to precipitate out or gas bubbles to form. Thus, there is a need for a gas cushion pressurization tank that maintains the single-phase state of a sample without requiring inordinately large and possibly dangerous pressures to be used in charging a sample tank before going down hole."}
{"text": "The present disclosure relates generally to content management services and in particular to selectively providing content management features within an existing messaging client.\nContent management services allow users to access and manage content across multiple devices using a network. In a typical content management service, a user establishes a content management account with the service provider and can associate various content items with the content management account. For example, content management services may allow the user to store content items (including but not limited to text documents; email messages; text messages; other types of messages; media files such as photos, videos, and audio files; and/or folders containing multiple files). Some content management services may allow users to set sharing options for stored content items in order to control whether and how other users are able to access the content items.\nMessaging services enable users to communicate with one another. The messaging services may include email, instant messaging (IM), or social networking services. In some cases, messaging services also allow for the transmittal of content items as attachments. For example, an email service may allow a user to compose an email, attach a file from the user's local machine, and send the file to a recipient through the email service. After the sent email is received, the recipient can open the email, download the attached file to the recipient's local machine, and view or otherwise use the attached file."}
{"text": "Both commercially and domestically, it is oftentimes desirable if not absolutely necessary to be able to locate studs beneath sheet rocked walls or as structural support members for roofs and the like. For even the simple task of hanging a picture or installing a new set of shelves, the finding of a structural stud member behind a surface can be important if not critical.\nBefore there were stud finders, one would either pound a small nail into the wall until it hit a stud or a small pivoting magnet would be used whose orientation would change when a magnetizable object such as a nail was confronted.\nElectronic stud finders replaced the hit and miss and pivotable magnet approaches. These products operate like a touch switch on a touch-sensitive lamp. They use changes in capacitance to sense stud location. When the plate inside the stud finder is over a surface such as wallboard, it will sense one dielectric constant but when over a stud, the dielectric constant is different. It works on a capacitance differential generated by density differences. The circuit in the stud finder can sense changes and report it on its display. The magnet approach was made the subject of U.S. Pat. No. 3,845,384 issued in 1974. The stud finder based upon detection of capacitive differences was made the subject of U.S. Pat. No. 4,099,118 issued in 1978.\nAlthough current commercially available stud finders operate reasonably well in the right conditions, there are many instances where the detection of capacitive differences behind a wall or membrane can be difficult if not impossible For example, in applying the hardware necessary to support photovoltaic panels on a roof, it is oftentimes necessary to locate structural stud members for applying standoffs which in turn support rail members and solar photovoltaic panels thereon. Installations of this type have been disclosed in applicant's previously filed U.S. application Ser. No. 12/407,952, the disclosure of which is incorporated by reference herein. However, roofs oftentimes support one or more layers of composite shingles and similar waterproof membranes which provide surfaces which are too thick and uneven to enable capacitive-based stud finders to function.\nIt is thus an object of the present invention to provide a new paradigm in stud finder technology which does not depend upon capacitive differences to enable a user to locate a structural stud member beneath a surface.\nIt is yet a further object of the present invention to provide a novel stud finder device which can be used in virtually any environment whether or not the studs to be located are sufficiently proximate a homogeneous surface to enable them to be detected by capacitive techniques.\nThese and further objects will be more readily apparent when considering the following disclosure and appended claims."}
{"text": "In recent years, businesses have found it more and more useful to analyze the transactions they undertake with their customers. In some cases, such an analysis can provide the business with the information it needs to fine tune the transactions to shape customer behavior to improve performance of the business.\nFor example, an electronic commerce company may want to analyze the “clickstream” of individual customers visiting the company's web page to find associations between customer actions, such as the customer clicking on a particular image on the web page, and desired customer actions, such as the customer making a purchase. With this information, the company may tune its web page to increase the likelihood that a customer will click on the image, with the hoped-for result that sales will increase. A method and apparatus for performing this analysis, called “association” or “affinity analysis,” using a massively parallel processing (MPP) computer system is described in co-pending U.S. patent application Ser. No. 09/410,528, entitled SQL-BASED ANALYTIC ALGORITHM FOR ASSOCIATION, filed on Oct. 1, 1999, and assigned to the same assignee as the present application.\nThis analytical approach can be useful in analyzing other forms of data, such as retail or financial data. The owners of a grocery store, for example, may find it useful in designing the layout of the store to know that customers are more likely to purchase cheese when they purchase grapes. Similarly, a bank may find it useful to know that customers who contract for certain banking services, such as a checking account, are more likely to acquire other banking services, such as automobile loans. The bank could use such information to design the mailing materials it includes with its statements.\nIt is sometimes important to isolate the actions of individual customers to properly perform affinity analysis. For example, in the web page context, in which a log is maintained containing customer transactions with the web page, it may be useful to associate transactions with individual customer “sessions” in order to allow a meaningful analysis of the transactions. A method and apparatus for performing such an analysis is described in co-pending U.S. patent application Ser. No. 09/752,355, entitled IDENTIFYING WEB-LOG DATA REPRESENTING A SINGLE USER SESSION, filed on Dec. 29, 2000, and assigned to the same assignee as the present application.\nUnderstanding the order of transactions may also be important. For example, a web page owner may be interested to know that a customer that clicks on a first image on the web page followed by a second image may be more likely to make a purchase than a customer that clicks on the second image before the first image. Making such a determination adds an extra degree of complexity to an affinity analysis."}
{"text": "Microelectromechanical systems (hereinafter MEMS) have been developed for movable micro devices, such as hygroscopes, accelerometers, tunable RF capacitors, digital mirrors, sensors and the like. They are used for forming small electrical and mechanical structures on a substrate, particularly a substrate of silicon or a silicon-containing material. These devices are made using conventional semiconductor processing techniques, such as chemical vapor deposition and plasma etching for example.\nFIG. 1 illustrates a conventional three-layer substrate that can be used to make a MEMS device. A substrate layer of silicon 10 is covered with a sacrificial layer of silicon oxide 11 and a layer of polysilicon 12 deposited thereover. In accordance with a prior art method, the silicon oxide layer 11 is etched away to at least partially separate the layer of polysilicon 12 from the substrate 10. This etch step is known as release. FIG. 2 illustrates a partial etch of the sacrificial silicon oxide layer 11, as by using an isotropic wet HF etch. Now the polysilicon feature 12 can move, e.g., bend toward and away from the silicon substrate 10. Suitable etchants are anhydrous and aqueous hydrogen fluoride (HF).\nThe above method however requires the steps of deposition and removal of the sacrificial layer, and requires a wet etch, as of HF, to etch through the sacrificial silicon oxide layer.\nIn addition to the extra steps required for depositing and isotropically etching away the sacrificial layer, release is also a problem. Release is a complex process wherein the silicon oxide layer must be controllably etched. If too much silicon oxide is removed, the desired structure is undercut; if too little silicon oxide is etched away, frozen microstructures are formed that are not able to move as intended. Further, such an etch produces residues which adhere to the substrate.\nAnother part of the problem is that different silicon oxides have different etch rates. Doped oxides, such as PSG, BPSG and doped TEOS oxide, attain a high and fairly stable etch rate quite rapidly. However, dense oxides, such as thermal oxides, TEOS oxides and high temperature oxides, have an etch initiation period, and a much slower etch rate, even after initiation. Etch rates can also vary with device design. Thus it is difficult to know how long the etch needs to be carried out for the desired release.\nAqueous HF has traditionally been used to isotropically etch silicon oxides. However, this etchant has a high surface tension that causes capillary and van der Waals forces that pull the released features toward the substrate until they contact each other. This results in a generally permanent adhesion of the feature to the substrate, a result known as stiction.\nAnhydrous HF has also been tried as a release etchant. However, it is a very strong acid, and will damage materials used in the interior of the etch chamber, adding to the costs of the process. A special chamber must be made, one that is at least partially impervious to anhydrous HF, a difficult and expensive challenge.\nThus a method of releasing MEMS devices that does not cause stiction, and that can be done in situ in a single chamber with fewer steps and high throughput that results in reduced costs, would be highly desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to an image recording apparatus in which is mounted a recording head that performs recording by ejecting (discharging) a liquid from an energy generating element or by thermal transfer.\nThe present invention can be applied for apparatuses, such as printers, copiers, facsimile machines for which communication systems are provided, or word processors that incorporate printers, that perform the recording of images on a recording medium, such as paper, thread, fiber, cloth, leather, metal, plastic, glass, wood or ceramics, and for industrial recording apparatuses with which various processors are combined.\nxe2x80x9cRecordingxe2x80x9d in this invention is defined not only as the formation on a recording medium of images, such as characters or drawings, that convey meaning, but also as the formation of images, such as patterns, that convey no meaning.\n2. Related Background Art\nConventionally, the demand for recording apparatuses that can produce high quality images has increased, and how to improve image quality has been the subject of numerous discussions. For a recording apparatus in which a recording head is moved in one direction when recording images, the precision of the positioning of an image to be recorded is determined by the accuracy with which the recording head itself is positioned. And for the improvement of the image quality, the enhancement of the accuracy with which a recording head is positioned is an extremely important element. Therefore, in a conventional recording apparatus, for a carriage on which is mounted a recording head that records in only one direction, position detection means (e.g., an image scanner) is provided for accurately ascertaining the position of the recording head. Or, at the carriage\"\"s home position in the apparatus, optical reading means is provided to detect the position of the recording head. Then, based-on the obtained head positioning data, whether the recording position is adequate or whether the recording position must be corrected is determined.\nHowever, in a conventional recording apparatus the recording head, which constitutes the printing means, and the position detection means are arranged separately. Therefore, in a recording apparatus wherein, for example, a head position detection means is provided for a carriage, satisfactory positioning accuracy for the recording head must be obtained by mounting the recording head on the carriage. In order to obtain such accuracy, precision in the sizing of components, such as the carriage and the recording head, must be improved, or a process must be performed for correcting the positioning of the recording head.\nIn addition, since elements and circuits for detecting the position of the recording head must be formed on the carriage or on the substrate of the apparatus, manufacturing costs will be increased.\nFrom the viewpoint of high quality image recording, highly delicate recording, for improved image density and tone representation, can be performed by producing dots that have variable sizes.\nAs the resolution of an image is increased, however, extremely high accuracy is needed to position the dots that are formed, and as the number of steps involved in varying the dot sizes is increased, greater dot size accuracy is required.\nThus, when a plurality of recording elements are employed, dot positioning errors and the use of non-uniform dot sizes can result in the deterioration of the image quality.\nIt is apparent that the demand for increased image quality can not be satisfied merely by improving the accuracy of the positioning of a carriage and a recording head and the accuracy in the production of dot sizes, so that accordingly, the shortcomings attributable to inaccurate dot positioning and to the unstable production of accurately sized dots are not resolved.\nIt is, therefore, one object of the present invention to provide at a low manufacturing cost an ink jet recording apparatus that can not only accurately detect the position of a recording head but can also accurately stabilize the positioning and the sizing of dots, a recording head therefor, and an element substrate to be used for the recording head.\nTo achieve the above object, according to one aspect of the present invention, provided is a recording head substrate on which are mounted energy generating elements that contribute to the formation of images by a recording head, and on which both light-receiving elements and light-emitting elements, or at least, light-receiving elements are mounted.\nThe light-receiving elements can be photodiodes or CCDs.\nIn addition, a controller for controlling the energy generating elements and the light-receiving elements is also mounted on the recording head substrate.\nIn this case, it is preferable that the light-receiving elements and at least one part of the controller be produced during the same manufacturing process.\nThe energy generating elements and the light-receiving elements are arranged along at least one line on the recording head substrate.\nThe energy generating elements and the light-receiving elements are arranged along a plurality of lines, and the lines are parallel to each other.\nIn this case, on the individual lines the number of the energy generating elements may be equal to the number of the light-receiving elements, but it is preferable that the number of the light-receiving elements be greater than the number of the energy generating elements.\nAccording to one more aspect of the present invention, provided is a recording head comprising:\nthe above described recording head substrate;\na top board in which are formed liquid flow paths that correspond to the energy generating elements; and\ndischarge orifices (port) which is communicated with the liquid flow path of the top plate and through which liquid is discharged by the application of energy by the energy generating elements,\nwherein the light-receiving elements and the light-emitting elements on the recording head substrate are optically opposite a face on which an image is formed by using the discharge ports.\nAccording to the present invention, as is described above the energy generating elements and the light-receiving elements are mounted on the same substrate. Therefore, when the light-receiving elements optically detect dots formed by the energy generating elements, accurate information concerning the positioning, the sizes and the densities of the image dots can be obtained quickly. Further, since in contrast to an arrangement where the energy generating elements, the light-emitting elements and the light-receiving elements are mounted on separate substrates, the process for the formation of the individual elements can be commonly employed and no connections are required, the manufacturing cost and the size of an apparatus can be considerably reduced.\nAccording to another aspect of the present invention, provided is a recording head substrate on which are mounted energy generating elements that contribute to the formation of images by a recording head, and on which are mounted a plurality of head position detecting elements for detecting the position of the recording head.\nAccording to an additional aspect of the present invention, a recording head, for forming images using energy generating elements, comprises:\na substrate on which are mounted not only the energy generating elements but also a head position detecting element for detecting the position of the recording head.\nAccording to a further aspect of the present invention, a liquid recording apparatus comprises:\na recording head for forming images employing energy generating elements while moving on a line;\nhead position detecting elements that are provided for the recording head for detecting the position of the recording head; and\na member in the recording apparatus that, in order to be detected by the head position detecting elements, is fixed opposite the head detecting element and along a track where the recording head moves.\nThe head position detecting elements are mounted on a substrate on which the energy generating elements are also mounted. In addition, it is preferable that, in accordance with position data for the recording head, detected by the head position detecting elements, and other recorded data, a circuit for generating signals to drive the energy generating elements, and light-receiving elements, for detecting an image that is formed, be mounted on the substrate on which the energy generating elements are mounted.\nThe head position detecting elements may be magnetic detecting elements, light-receiving elements or electric field detecting elements. The energy generating elements may be electro-thermal converting elements for heating liquid and inducing film boiling in order to discharge liquid droplets for forming images.\nAs is described above, according to the present invention, since the energy generating elements that contribute to image recording and the elements for detecting the position of the recording head are mounted on the same substrate, the accuracy at which the position of an image can be recorded is extremely high. In addition, since using semiconductor fabrication processing at least the elements having two functions can be mounted on the same substrate at the same time, the manufacturing costs can be drastically reduced."}
{"text": "The present invention relates to a stereo image photographing system, and more particularly to one capable of easily obtaining information and measuring data necessary for obtaining three-dimensional data from a stereo image.\nIn recent years, line drawing has been made for a site to be measured by using a surveying instrument, e.g., a pen-based mapping system, and a portable computer in combination thereof. In such a conventional pen-based mapping system, plotting has been carried out by drawing a line image on the portable computer, using a surveying instrument such as a total station and a GPS. The plotting of the site to be measured was carried out on site by watching the situation there, or in an office based on photographed images. In aerial photography, plotting was carried out by using an analytical plotting device or the like based on stereo-photographed images. When three-dimensional measurement was carried out from a stereo image in a terrestrial photogrammetry, conventionally, photographing, analysis and measurement were executed by a stereo-camera having a base line length and the facing direction of the camera strictly fixed.\nHowever, in the conventional art, the fixing of the base line length or the like of the stereo-camera has placed limitations on the range of photographing and accuracy (object). In addition, conventionally, when measurement was carried out without fixing the base line length or the like, even if two images were overlapped and photographed for three-dimensional analysis by stereo-photographing, it was difficult to form an image to be subjected to stereo-analysis. In other words, in most cases, the analysis of images was impossible, and even if possible, the images were unstable with low accuracy and considerably low reliability. In particular, when photographing was carried out with a camera loaded on a balloon, a car, or other moving body, it was extremely difficult to obtain a highly reliable and analyzable image.\nThe present invention was made with the foregoing problems in mind, and it is an object of the invention to easily obtain a stable, highly reliable and accurate stereo image when three-dimensional measurement is carried out by a general camera. It is another object of the invention to provide a stereo image photographing system which makes it easy and reliable to perform stereo-photographing for three-dimensional measurement, which has been difficult, and further enable automatic stereo-photographing.\nIn accordance with a first aspect of the present invention, a stereo image photographing system is provided, comprising:\na photographing unit which is movable and has a camera for photographing and a reflection member placed thereon;\na signal generating unit for generating a timing signal indicating a photographing timing when photographing is made by the camera of the photographing unit; and\na measuring unit placed in a position apart from the photographing unit, for emitting a measuring light toward the reflection member of the photographing unit in accordance with the timing signal from the signal generating unit, and measuring a position of the photographing unit based on a reflected light reflected from the reflection member.\nIn accordance with a second aspect of the present invention, a stereo image photographing system is provided, comprising:\na photographing unit which is movable and has a camera for photographing an image of an object in accordance with a timing signal indicating a timing and a reflection member placed thereon;\na measuring unit placed in a position apart from the photographing unit, for emitting a measuring light toward the reflection member of the photographing unit, and measuring a position of the photographing unit based on a reflected light reflected from the reflection member;\na photographing position calculating unit for calculating photographing position data based on photographing condition data regarding the object to be photographed;\na driving unit for moving the photographing unit to a predetermined position based on the photographing position data calculated by the photographing position calculating unit;\na movement control unit for controlling the driving unit; and\na photographing control unit for outputting a timing signal to the camera of the photographing unit when coincidence/rough coincidence is determined between position data regarding the position of the photographing unit measured by the measuring unit and the photographing position data calculated by the photographing position calculating unit.\nIn accordance with a third aspect of the invention, a stereo image photographing system is provided, comprising:\na photographing unit which is movable and has a camera for photographing and a reflection member placed thereon;\na photographing position calculating unit for calculating photographing position data based on photographing condition data regarding an object to be photographed;\na driving unit for moving the photographing unit to a predetermined position based on the photographing position data calculated by the photographing position calculating unit;\na movement control unit for controlling the driving unit; and\na storing unit for storing data of an image photographed by the camera of the photographing unit."}
{"text": "Natural coal is generally regarded as a non-renewable energy source in part because of the amount of time needed for the formation thereof to occur naturally. For instance, lignite coal requires millions of years, and anthracite coal requires hundreds of millions of years for its respective formation. The burning of natural coal releases toxic agents, such as heavy metals (e.g., mercury, uranium) co-deposited as part of its sedimentary formation process. Other serious toxic by-products resulting from the burning of coal and other fossil fuels such as oil, gasoline, tar and natural gas include: soot, ash, and oxides of nitrogen, sulfur, and carbon. The release of these toxic chemicals into our atmosphere from the burning of these fossil fuels for the production of energy has resulted in air pollution as well as global warming. Additionally, since fossil fuels are a finite resource and the recovery thereof is becoming more costly as the supplies thereof dwindle, continued reliance thereon will at some point become economically as well as environmentally unfeasible. Also, since a significant amount of crude oil used in the United States is imported, a substantial imbalance in the balance of trade with oil-exporting countries has resulted in severe economic damage to US economy due to the outflow over time of trillions of dollars.\nA further environmental problem concerns the need for nitrogen fertilizers for agricultural purposes. Currently, nitrogen fertilizers such as ammonium nitrate and urea require ammonia for their manufacture. Ammonia is chiefly manufactured from nitrogen and hydrogen at high temperature and pressure employing the Haber-Bosch process. While nitrogen can be obtained relatively easily from the atmosphere, hydrogen is obtained by a process known as steam reforming from natural gas. Currently, 5% of the world's annual natural gas is utilized for that purpose. Thus, it would be desirable to generate nitrogen fertilizers without reliance on the use of natural gas.\nIn response to significantly increased crude oil prices and politically and economically uncomfortable dependence on imports of crude oil, ethanol has been utilized as an additive to gasoline since 1980 and has now commonly reached concentrations of 10% by volume, with 20% by volume being legislated to occur in Minnesota in 2013. Several states also currently offer E85 which contains 85% ethanol and 15% gasoline. A primary advantage of synthetic fuels derived from biomass is that the carbon in them that is released as carbon dioxide into the atmosphere when burned is not newly introduced and additive to the carbon dioxide already present therein. Rather, the released carbon dioxide is recaptured and recycled by plants through photosynthesis resulting in no net addition of that global warming gas to the earth's atmosphere.\nThe prior art process of producing ethanol and which is representative of how 85% of the ethanol produced in the US is made, can be understood by referring to FIG. 3. As seen therein ethanol, a two carbon alcohol and, to a lesser degree, butanol, a four carbon alcohol, are obtained in the United States primarily by the fermentation in large fermentation vats 10 containing a ground fermentable material. Typically, such material is starch as present in corn kernels, potatoes and other like vegetable materials. Water is added to the fermentable material and enzymes, natural and/or added, are used to breakdown the starch into fermentable dextrins and sugars. Of course, it is also well known to simply start with sugars as obtained, for example, from sugarcane or sugar beets. Yeast is also added to the vat 10 and the mixture is heated by a heater 12 to approximately 37 degrees C. to facilitate the fermentation of the dextrins and sugars producing the ethanol and generating carbon dioxide gas. At the end of the fermentation process the fermented material in vat 10 is pumped along line 14 to a “beer well” 16 which also includes a heater 18. Well 16 is a sealable container and a distillation apparatus 20 is fluidly secured thereto by an outlet conduit 22. The fermented material is then heated by operation of heater 18 to a temperature sufficient to boil the ethanol but not the water, i.e. between approximately 78 and 90 degrees C. As is well understood, the ethanol vapor is then cooled and condensed to liquid form by distillation apparatus 20 and directed to a storage container 24. As is well understood, the resultant ethanol product will contain small amounts of water because water and ethanol form an azeotropic mixture that cannot be separated by distillation. As is also well understood, subsequent processes can be applied such as through the use of molecular sieves for the conversion of the distilled ethanol into absolute or 100% ethanol. After the removal of the ethanol from the fermented material, what remains is a fermentation residue. Those of skill will recognize that the fermentation residue is a highly processed material which differs from the original ground fermentable materials in it is depleted of carbon both from the loss of the carbon dioxide gas during the fermentation step and from the ground fermentable material and the fermented residue as it is further depleted of carbon through the loss of the carbon contained in the ethanol or butanol that is removed from the fermented residue by distillation.\nAfter the removal of the ethanol, the remaining fermentation residue in well 16, referred to as whole stillage (WS), having a dry solids content of about 10-15%, in the dry-grind corn ethanol production industry, is then sent to a centrifuge 26 where it is separated into a solid portion referred to as Wet Distillers Grains (WDG), having a dry solids content of about 35%, and a liquid portion called Thin Stillage (TS), having a dry solids content of 5-10%. TS can be further processed by the evaporation of the water therefrom in evaporator 28 to produce a thick syrup-like portion referred to as Condensed Distillers Solubles (CDS), having a dry solids content of 30-50%. WDG can be sent along a conduit depicted by dashed line 32 directly to a dryer 30 and dried into a shelf-stable animal feed product called Dried Distillers Grains (DDG) by operation of a heater 34 thereof. Alternatively, a shelf-stable animal feed product referred to as Dried Distillers Grains with Solubles (DDGS) can be manufactured. DDGS is typically what is made in the ethanol industry and is obtained by the combining of WDG and CDS, to form a mixed grains material followed by extensive drying of that combination in dryer 30 until the DDGS has a moisture content of less than 10% dry weight and is used as a shelf-stable animal feed product. Although DDG is an acceptable, high quality shelf stable animal feed as well, the DDGS includes additional nutrients present in the CDS. Typically, a bushel of corn can provide 2.8 gallons of ethanol, 18 pounds of DDGS, and 16 pounds of carbon dioxide.\nThe WDG fermentation residue is comprised of residual fermentable dextrins, yeast, non-fermentable dextrins, protein, fat, low concentrations of metal salts, and water. Often the protein content can be quite high, ca. 30%, and the material is utilized as a high-grade feed especially for ruminants such as cattle, with small amounts able to be employed in feeds for mono-gastric animals such as swine and poultry. WDG can be transported directly to farm sites but only within a relatively small radius, e.g., 30 miles, of the manufacturing site. The limited shipping radius of 30 miles is due to the amount of water present in the WDG which leads to microbial growth that quickly renders it unsuitable for animal consumption.\nUnfortunately, and in either case, the physics of the removal of water by heating in dryer 30 requires the input of substantial amounts of energy which adds considerable cost to the production of either DDG or DDGS feed. It is estimated that the intensive drying operations performed on these by-products accounts for 40-45% of the thermal energy and 30-40% of the electrical energy required in the dry-grind ethanol process. In addition, there is also a significant energy cost associated with evaporating the TS and producing CDS. To complicate matters, ethanol manufacturers are often in the situation of having a large surplus of WDG that must be dried at the manufacturing site as there often exists insufficient need for it within the 30 mile radius of their plants. Other fermentation residue stillages that can have similar negative economic and/or waste disposal issues include Brewer's Grains from the production of beer and whiskey, E.-Coli residues from production of antibiotics, and residues from other bacterial, fungal and yeast fermentation processes. As a result, there has been a long felt need for decades to find an economical, energy-efficient and environmentally sustainable solution for obviating the need to dry fermentation residues such as TS and WDG.\nThis need is particularly pressing for the production of biofuels because it is has been difficult, if not impossible, for the alternative fuel industry to compete on a dollar per gallon basis with the petroleum industry. And specifically, the ethanol industry has been criticized for not being energy efficient, that is, the energy represented by the ethanol being barely greater or even less than the energy needed to produce it, especially when the drying of the TS and WDG are taken into account. However, despite the tremendous pressure to find solutions that would create an economically viable and sustainable alternative fuel industry, results have not been forthcoming.\nThere also exists a need in the energy resource field for improved synthetic coal products that possess a commercially viable high British Thermal Unit (BTU) per pound content. There further exists a need for nutritive media that can be utilized as plant fertilizers and reduce the consumption of natural gas required in the manufacture of nitrogen fertilizers.\nAdditionally, there exists a need for innovative processes utilizing natural biological materials as carbon neutral energy resources, thereby not increasing the net amount of carbon dioxide in the atmosphere that utilizing fossil fuels causes."}
{"text": "The invention relates to a locking device of a cover for closing an opening in a wall of a storing compartment, in particular in a trunk of an automotive vehicle.\nGenerally such a locking device includes a pivotal handle for displacing a locking bar from a locking position to an unlocking position against the force of a return spring. When the locking bar is in the locking position, it retains the cover in its position within the trunk, and when the locking bar is in the unlocking position, it releases the loading cover so that the cover may be lifted from the opening.\nThere is a certain risk that the locking bar, when the cover is lifted from and inserted into the opening, may move to its unlocking position where it projects relatively far from the locking device. As a result the edge of the opening within the trunk, which is generally covered by a textile layer, may be damaged. Furthermore, the projecting locking bar renders lifting and inserting the loading cover from and into the opening relatively difficult."}
{"text": "1. Field of the Invention\nThe document relates to a system and a method for controlling transport rates of real time streaming services. The services are configured to control a transport rate of data transmitted from a streaming server to a client using a real time transport protocol, thereby maintaining a buffer level of a streaming client at a constant level.\n2. Discussion of the Related Art\nConventionally, multimedia data, such as video data, computer animation data or any combination thereof, is all downloaded to a hard disc drive of a PC (Personal Computer), for example, via access to an Internet for reproduction of the same. However, in the streaming environment, the entire data may not be all downloaded, and the data can be reproduced real time in the course of buffering without first downloading the entire file data. In other words, streaming is the ability to view data without first downloading an entire file. In general, streaming is a multimedia data transmission method which reads, transmits and plays data on an Internet in real time.\nWhen a streaming server and a streaming client are connected over a wired/wireless network, and the streaming server transmits the multimedia contents, the streaming client having received the transmission of the multimedia contents buffers the received contents and reproduces the contents real time.\nBasically, the transport rate of the streaming server conforms itself to requirements of data to be transmitted. In audio/video stream packets transmitted by compression via Moving Picture Experts Group (MPEG), time information for controlling the transport rate of the streaming server is recorded in a header.\nHowever, there arises a problem even if the transport rate is under prolonged time of streaming service. It is impossible to calculate an accurate transport speed even if the time information is of an extreme accuracy, and a small difference of transport speeds can cause a buffer overflow after prolonged time.\nThere is another problem in that the streaming server and the streaming client, each being an independent identity, have separate reference clocks, and fine differences of clocks are accumulated in the course of use of a cheap crystal oscillator to prevent provision of a stable operation.\nFIG. 1 is a first graphical representation illustrating a buffer level change according to the related art.\nAs in FIG. 1, the buffer level keeps increasing if the transport speed of the streaming server is faster than the speed of data reproduced by an actual streaming client, resulting in a case of exceeding a maximum buffer level defined by the streaming client.\nAnother problem is that a method of determining a transport rate only by time information contained in the multimedia data to be transmitted causes to generate a time delay from a streaming start time to a start of an actual reproduction.\nFIG. 2 is a second graphical representation illustrating a buffer level change according to the related art where a buffer level change following the start of the streaming is shown. There is needed a minimum quantity of data for grasping and reproducing context of multimedia content in which an MPEG decoder of a streaming client is transmitted. This is because the streaming client can accept a delay variance to a certain degree. Subsequently, as in FIG. 2, ‘a buffer level to be required’ supposed to be filled in the buffer without reproducing the transmitted data is determined up to a certain level prior to the start of the reproduction.\nHowever, control by a simple comparison of buffer level reported from the required buffer level and the streaming client can deteriorate the stability of the system, and a very small difference of transport speeds between the streaming server and the streaming client can result in generation of a buffer overflow or an under-run.\nThe streaming server is required to continuously control a transport rate in order to maintain a buffer level reported by the streaming client. Furthermore, there is a need of filling the buffer at a maximally fastest speed at a point where streaming starts for a fast reproduction."}
{"text": "1. Field of the Invention\nThe present invention relates to a watertight connector.\n2. Description of the Related Art\nA known watertight connector has a housing with opposite front and rear ends. A fittable portion is disposed at the front end of the housing and is dimensioned to fit in a receptacle of a mating connector. Terminal fittings are mounted in the fittable portion and a tubular seal is mounted on the outer peripheral surface of a fittable portion. The seal seals a clearance between the outer peripheral surface of the fittable portion and the inner peripheral surface of the receptacle.\nA holder is provided at the rear end of the fittable portion of the housing and a lock is formed at the rear end of the seal to prevent the seal from coming off the fittable portion when the watertight connector is withdrawn from the receptacle. The lock engages the holder to prevent loose forward movements of the seal and separation of the seal from the fittable portion.\nIt generally is desirable to shorten the seal along forward and backward directions to reduce connection resistance with the receptacle and to prevent the seal from buckling. Thus, the seal is at a rear end of the fittable portion and is held resiliently in contact with a front area of the receptacle.\nThe mating connector often is on a fixed piece of equipment and the watertight connector may have to be fit into the receptacle by touch due to a narrow space near the receptacle. Thus, the opening to the receptacle may be slanted outwardly to guide the insertion of the fittable portion. However, the degree of resilient deformation of the seal with the slanted surface changes depending on the depth of engagement of the fittable portion, and sealing performance is unstable. Accordingly, the slanted leading end cannot be a sealing surface. A portion of the receptacle more backward than the slanted opening is set as a sealing surface for resilient contact with the seal, and a resilient contact of the seal with the receptacle is near the front-end of the fittable portion.\nThe rear end of the seal is held at the rear end of the fittable portion so as not to come off. Thus, the seal must be longer along forward and backward directions if an attempt is made to position the resilient contact of the seal at the front end of the fittable portion. However, frictional resistance pushes the resilient contact portion of a longer seal back during the fitting operation, and hence the longer seal may buckle.\nThe present invention was developed in view of the above problems and an object thereof is to prevent a sealing member from being buckled.\nThe invention relates to a watertight connector with a housing in which at least one terminal fitting is accommodated. The watertight connector has a fittable portion that can be fit in and withdrawn from a receptacle of a mating connector. A substantially tubular seal is mounted on the outer peripheral surface of the fittable portion and is formed on its outer peripheral surface with at least one resilient contact that is held resiliently in contact with the inner peripheral surface of the receptacle. One or more fastening means are provided at the rear end of the fittable portion and/or at the rear end of the seal for preventing the seal from moving forward and off the fittable portion. One or more loose-movement preventing means are provided on the fittable portion and on the seal for preventing loose backward movement of the resilient contact.\nThe fastening means preferably are provided at the rear end of the fittable portion and/or at the rear end of the seal.\nThe loose-movement preventing means preferably is at the fittable portion and at the seal and engage each other to prevent loose backward movement of the resilient contact.\nFrictional resistance that acts during the fitting of the watertight connector into the receptacle may push the resilient contact back. However, the loose-movement preventing means prevents the resilient contact from moving loosely backward in response to forces exerted by the receptacle. Thus, the seal will not buckle even if the seal is long along forward and backward directions and the resilient contact is at the front-end of the seal.\nThe loose-movement preventing means preferably comprises steps on the outer peripheral surface of the fittable portion and on the inner peripheral surface of the seal. The steps prevent the loose movements of the resilient contact by bringing into contact surfaces of the steps aligned at an angle to a fitting direction of the fittable portion into the receptacle.\nA backward-acting pushing force exerted by the receptacle on the resilient contact is received by the step of the fittable portion via the step of the seal. The surfaces of the steps intersect the fitting direction of the fittable portion into the receptacle. Thus, the fittable portion can receive the pushing force from the receptacle to prevent loose backward movements of the resilient contact.\nThe steps of the fittable portion and the seal prevent the loose movements of the resilient contact by bringing into contact surfaces that overhang or are undercut with respect to the fitting direction. Thus, an acute-angled corner of each step bites in or engages an angled corner of the mating step due to the inclination of the overhanging surface to prevent the loose movements of the resilient contact. Accordingly, disengagement of the steps and loose movements of the resilient contact can be avoided.\nThe step of the seal is at a position near the resilient contact with respect to forward and backward directions.\nThe backward pushing force of the receptacle on the resilient contact deforms the resilient contact inwardly and the restoring force of the resilient contact presses the step of the seal against the outer peripheral surface of the fittable portion. Accordingly, the step of the seal will not disengage outwardly from the step of the fittable portion and loose movements of the resilient contact are prevented.\nAn escaping groove is formed in an area of the inner surface of the fittable portion that faces the terminal insertion space in the housing. The escaping groove avoids interference with an improper insertion preventing projection on the outer surface of the terminal fitting when the terminal fitting is inserted in a proper orientation. An area of the outer surface of the fittable portion corresponding to the escaping groove preferably forms a rib shape, and the front end of the rib-shaped projection defines the step.\nA nonuniform thickness of the wall of the fittable portion can be reduced by forming the rib shape on the outer surface of the fittable portion to align with the escaping groove. Thus, a variation in the degree of deformation caused by thermal shrinkage (i.e. xe2x80x9csink marksxe2x80x9d) during the resin molding of the housing can be suppressed to ensure dimensional stability. Further, the step of the fittable portion takes advantage of the rib-shaped projection corresponding to the escaping groove. Accordingly, the shape of the fittable portion can be simplified as compared to a case where the step is separate from the rib-shaped projection.\nThe fastening means may comprise a stopper fit to the housing for engaging the seal and preventing the seal from moving forward.\nA single loose-movement preventing means may be provided on each of the fittable portion and the seal around substantially the entire periphery thereof.\nAlternatively, a plurality of loose-movement preventing means may be spaced around the periphery of the fittable portion and/or the seal.\nThese and other objects, features and advantages of the present invention will become more apparent upon reading of the following detailed description of preferred embodiments and accompanying drawings. It should be understood that even though embodiments are separately described, single features thereof may be combined to additional embodiments."}
{"text": "Conventional portable audio systems often include a pair of headphones that are connected to a media player (e.g., by one or more wires or by wireless technology). It is increasingly common for users to use portable audio systems when engaging in outdoor activities. While the media player in any given portable audio system can be used in a variety of settings, it is often the case that the headphones employed are not as versatile. For example, in-ear headphones may provide for portability, but such headphones may provide poor audio quality, be uncomfortable, or both. Where multiple wires are used to connect the headphone speakers to the media player, each additional connection (e.g., each connection between male and female audio jacks) may further degrade audio quality. While larger, over-the-ear headphones may be more comfortable, they may be awkward to wear with outdoor gear, such as goggles. For example, it is increasingly common for outdoor enthusiasts, such as skiers and snowboarders, to use portable audio systems when engaging in outdoor activities, such as skiing and snowboarding. In most cases, skiers and snowboarders favor smaller, in-ear style headphones because helmets, ski goggles, ear protectors, hoods, and headbands can more easily be worn over such headphones."}
{"text": "This invention relates generally to hinge means and, more particularly, to hinge means for mounting an element such as a door for swinging between open and closed positions on an element such as a cabinet.\nIn certain so-called European-style cabinets, the door is mounted on the cabinet by hinges which are fully concealed from the front and sides of the cabinet when the door is fully closed. Prior fully concealed hinges of this type are comparatively large and complex and require a significant amount of material. For example, a hinge which enables the door to be opened through approximately 175 degrees generally requires a linkage mechanism with as many as seven parallel pivots or requires a linkage mechanism with six parallel pivots and a slide. Another type of a prior fully concealed hinge has a somewhat simplified linkage mechanism with only four parallel pivots but such a hinge generally does not permit the door to be opened beyond approximately 125 degrees.\nThe complexity of the prior fully concealed hinges is increased even further when the hinges are of the self-closing type. A self-closing hinge includes a spring mechanism which acts to swing the door to its fully closed position once the door has been closed part way."}
{"text": "Receivers that detect stereophonic/monophonic signals are incorporated into a vast number of devices used in everyday life. For example, such receivers are used in automobile radios, a variety of communication systems such as cellular telephones, and even in children's toys. Unfortunately, many modern receiver systems suffer from performance shortfalls, such as frequent switchover between monophonic and stereophonic modes due to noisy channel environments and false detection of stereophonic signals as monophonic due to rated maximum system deviation (RMSD) mismatch.\nIn order to receive FM audio signals, be they music or any other type of information, a receiver must be robust enough to handle changes in the channel wherein the transmission could become very noisy and/or must overcome interference. Generally, a pilot tone is transmitted as part of the baseband signal that is used to modulate an FM carrier signal in order to indicate the nature of the transmission to be stereophonic. The energy of the pilot tone may fluctuate significantly in a harsh channel scenario. Simply comparing the pilot tone energy, estimated at the receiver, against a predetermined threshold may cause the receiver to switch between monophonic and stereophonic mode too frequently and degrade the entertainment quality of the audio program delivered to the consumer.\nIn addition, the receiver structure and the accompanying algorithms must also be flexible enough to handle a situation where the transmitted FM signal RMSD is not known beforehand. Generally, the allowed RMSD values are 75 kHz and 50 kHz. Hence, a mono/stereo transmission may be utilizing either one of them. However, the receiver may be set to operate at a different RMSD than the received signal RMSD. If the received signal RMSD and the RMSD of the receiver are not matched, a situation may occur wherein a stereophonic signal may be falsely interpreted as monophonic by the receiver. This results in the listener being denied the stereophonic quality of the program that the service provider is transmitting on the airwaves."}
{"text": "The present invention relates generally to microfabrication methods. More particularly, the invention relates to methods for forming and positioning moldable permanent magnets on electromagnetically actuated microfabricated components, such as electrical switches. Micromachining is a recent technology for fabricating micromechanical moving structures. In general, semiconductor batch fabrication techniques are employed to achieve what is in effect three-dimensional machining of single-crystal and polycrystalline silicon and silicon dielectrics, producing such structures as micromotors, microsensors and switches. Thus, except for selective deposition and removal of materials on a substrate, conventional assembly operations are not involved.\nElectromagnetically actuated micromechanical components include permanent magnets attached to a microfabricated moving structure. A separate electromagnet, appropriately energized, produces a controllable magnetic field which cooperates with the permanent magnet to cause the microfabricated element to move.\nOne example of an electromagnetically actuated microfabricated component is disclosed in an article by B. Wagner and W. Benecke, \"Microfabricated Actuator with Moving Permanent Magnet\" IEEE MEMS '91, Nara, Japan (1991). Another example is disclosed in commonly-assigned Ghezzo et al. U.S. patent application Ser. No. 08/000,172, filed Jan. 4, 1993, entitled \"Micromechanical Moving Structures Including Multiple Contact Switching System, and Micromachining Methods Therefor.\"\nA temporary structure, known alternatively as a release layer or as a sacrificial layer, is critical for micromachining because it allows moving parts to be formed by self-registered casting methods, with subsequent selective etching to remove the sacrificial layer. Since micromachining originated from the technology of silicon integrated circuit processing, low-temperature SiO.sub.2 is commonly employed as a sacrificial layer material. However, metals may also be employed for the sacrificial layer. For example, the use of a copper sacrificial layer is disclosed in the above-identified application Ser. No. 08/000,172.\nThe fabrication of such structures also requires that magnets somehow be attached to or formed on a movable element. In general, two approaches have previously been employed.\nThe first approach has been to attach magnets using pick-and-place equipment. This approach, however, has several disadvantages. In particular, the overall process is relatively time-consuming, in that suitable magnets must be purchased, and the magnets must be cut to a specific size, which is difficult with tiny elements. The magnets are then carefully placed and secured in position employing a suitable adhesive. Thus, while rare earth NdFeB magnets are effectively used in miniature components as small as 1 mm, they are difficult to machine, and in general are costly to produce.\nThe second approach has been to deposit magnetic material employing RF sputtering techniques. It has been observed, however, that crack-free rare earth transition metal films cannot be deposited with thicknesses greater than about 10 .mu.m. Such is reported for example in F. Cadieu, \"Permanent Magnet Thin Films: A Review of Film Synthesis and Properties\" Physics of Thin Films, M. Francombe and J. Vossen, Editors, Volume 16, Academic Press Inc., 1992. In addition, to obtain the desired physical properties, the films need to be heated to temperatures greater than 400.degree. C. While such may be possible in the case of silicon micromachining in general, such temperatures cannot be used in the presence of materials employed in high density interconnect (HDI) technology.\nMore particularly, what is known as high density interconnect (HDI) technology for multi-chip module packaging, is disclosed for example in Eichelberger et al. U.S. Pat. No. 4,783,695. Very briefly, in systems employing this high density interconnect structure, various components, such as semiconductor integrated circuit chips, are placed within cavities formed in a ceramic substrate. A multi-layer overcoat structure is then built up to electrically interconnect the components into an actual functioning system. To begin the multi-layer overcoat structure, a polyimide dielectric film, such as Kapton polyimide (available from E. I. DuPont de Nemours & Company, Wilmington, Del.), about 0.5 to 3 mils (12.7 to 76 microns) thick, is laminated across the top of the chips, other components and the substrate, employing Ultem.RTM. polyetherimide resin (available from General Electric Company, Pittsfield, Mass.) or another thermoplastic as an adhesive. The actual as-placed locations of the various components and contact pads thereon are determined by optical sighting, and via holes that are adaptively laser drilled in the Kapton film and adhesive layers in alignment with the contact pads on the electronic components. Exemplary laser drilling techniques are disclosed in Eichelberger et al. Pat. Nos. 4,714,516 and 4,894,115; and in Loughran et al. Pat. No. 4,764,485, each of which is incorporated by reference. Such HDI vias are typically on the order of one to two mils (25 to 50 microns) in diameter. A metallization layer is deposited over the Kapton film layer and extends into the via holes to make electrical contact to chip contact pads. This metallization layer may be patterned to form individual conductors during the deposition process, or it may be deposited as a continuous layer and then patterned using photoresist and etching. The photoresist is preferably exposed using a laser which is scanned relative to the substrate to provide an accurately aligned conductor pattern upon completion of the process. Exemplary techniques for patterning the metallization layer are disclosed in Wojnarowski et al. U.S. Pat. Nos. 4,780,177 and 4,842,677; and in Eichelberger et al. Pat. No. 4,835,704 which discloses an \"Adaptive Lithography System to Provide High Density Interconnect\" each of which is incorporated by reference. Any misposition of the individual electronic components and their contact pads is compensated for by an adaptive laser lithography system as disclosed in U.S. Pat. No. 4,835,704. Additional dielectric and metallization layers are provided as required in order to make all of the desired electrical connections among the chips.\nAs is described in detail hereinbelow, one aspect of the present invention is the fabrication of molded permanent magnets on micromechanical structures employing a modified HDI technology, including use of a polymer such as a polyimide (e.g. Kapton) as a mold material."}
{"text": "This invention concerns axial fluid flow inducing devices such as fans, which use a set of rotating blades reacting with the fluid to cause axial fluid flow.\nIn the above-referenced copending application there is described a unique direct magnetic drive in which impeller blades are permanently magnetized to establish a magnetic pole at each blade tip, which interacts with pulsed magnetic fields set up by one or more electromagnets located around the perimeter of the impeller or rotor in such a way that magnetic repulsion and attraction forces acting on the blades cause the impeller to rotate.\nIn the past, in order to increase the static pressure generated by an axial flow impeller, it has been the practice to provide fixed stator blades providing fluid reaction surfaces increasing the static pressure generated downstream of the impeller. The presence of the stator blades creates flow resistance, reducing flow. To avoid this disadvantage, it is known that counter-rotating impellers increase the static pressure without increasing flow losses in the system. However, complex mechanical drives are necessary in order to drive two or more closely positioned coaxial impellers in opposite directions.\nIt is one object of the present invention to provide a drive for multi counter rotating impellers not requiring a complex mechanical drive means.\nAnother known flow device has multiple impellers connected together and rotating in the same direction in order to provide a multiple stage pumping action. However, in order to enable proper multistage operation of synchronously rotating impellers, reduced blade areas must be established that increased the cost. In this application blade area remains the same while rotation speed is changed.\nIt is thus another object of the present invention to provide a multistage impeller in which the material and labor costs are reduced.\nThe above objects as well as others that will become apparent upon reading of the following specification and claims are achieved by an arrangement of a plurality of impellers that are coaxially mounted in a shroud and independently rotated with respect to each other. In a first application of the invention, the impellers are driven by forces created by intermittently generated magnetic fields acting on magnetized portions of each of the impellers so as to rotate the impellers. By causing rotation of two impellers in opposite directions a much-simplified means is provided for obtaining increased static pressure of the fluid flow induced by rotation of the impeller blades. Preferably, the magnetic field is generated by a pulsed generating means including one or more driver coils each having an U-shaped core with opposite core legs extending generally radially and with their ends located adjacent the outer perimeter of the associated impeller.\nPreferably, a pair of electromagnetic driver coils is provided which are arranged at an angle to each other and spaced apart so that when one pair of core legs straddles one blade, the other pair has sequential blades aligned with each leg. A pair of sensors, such as Hall effect sensors, are arranged to detect the passage of the leading edge of each successive impeller blade and to control the energization and magnetic polarity of the driver coils such as to induce rotation in opposite directions of each impeller by magnetic interaction between the field of the respective coils and the magnetic field of the permanently magnetized portions of the impeller blades.\nFor the embodiment of the invention having two counter rotating impellers, each electromagnetic coil is preferably skewed such that respective legs of the coil core are located at the outer perimeter of a respective impeller so that both impellers are energized and driven by the same pair of electromagnetic coils making synchronization of the rotation of the two impellers much simpler. However, using more than one pair of electromagnetic coils is also possible, each disposed around the perimeter of a respective impeller.\nDriving the impellers by a single driver coil is also possible. In this case, a simple electromagnet may be utilized in order to locate each impeller in a proper start up position with respect to the driving electromagnetic coil in which the spaced core legs straddle a blade tip.\nThe impellers preferably have blades of a type of plastic material that is permanently magnetizable and the magnetized tips comprise the outer portions of the impeller blades that are magnetized to interact with the pulsed magnetic fields.\nThe impellers may be constructed in two sections, each having an alternate set of impeller blades with the tips magnetized with the same polarity, and each impeller section having its tips magnetized with opposite poles to the other section. The impeller sections are interfit at assembly to produce an impeller in which successive blades are magnetized with opposite polarities.\nAlternatively, instead of a stationary shroud, the fluid flow passage can be formed by an outer ring that concludes, which is formed by a series of arcuate segments attached to the tip of each rotor blade, each segment interfit to the next adjacent segment on either side. Again, a two-section impeller construction may be advantageously employed.\nAn intermediate ring fixed to one impeller may also be employed with a labyrinth seal formed between the ring and adjacent end of the corresponding impeller.\nThe impeller ring segments adjacent the housing end wall may be provided with projecting portions that are received within the housing to form a labyrinth seal, such that while each ring rotates independently of the other, a sealed confinement of the fluid flow is assured.\nIn another version, the outer ring can be attached to the tips of blades and magnetized by segments with opposite polarities.\nThe impellers may also have magnetized outer portions which are driven by pairs of driver coils in such a manner as to be driven in the same direction at controllably different rates of rotation such as to produce a multistage pumping action which does not require differing impeller blade areas. A set of stator blades is required to be placed between impellers.\nIn another version, connected impellers can be disposed on either side of a set of stator blades."}
{"text": "1. Field of the Invention\nThe present invention relates to a silicon carbide (SiC) structure, a method of forming the SiC structure, in particular, the invention relates to a single crystal SiC structure, a method of the single crystal SiC structure, and an electronic device having the single crystal SiC structure.\n2. Background Arts\nAn electronic device implementing an SiC substrate has been drawn attention because of not only excellent breakdown voltage but also higher operable frequencies thereof. For instance, an electronic device operable at higher temperatures, having excellent breakdown voltages and losses in higher frequencies is realizable by using a channel made of silicon carbide (SiC) like an SiC-MOSFET (SiC metal-oxide-semiconductor filed effect transistor). Also, an electronic device with a graphene layer on the SiC substrate may be operable at higher frequencies.\nA Japanese Patent Application laid open No. of JP-2004-022878A has disclosed an SiC-MOSFET having (000-1) crystal plane as a primary surface thereof in order to enhance carrier mobility in a channel. An international Patent Applications laid open No. WO2009/063844 has disclosed an SiC substrate with an non-polar surface macroscopically but with a composite surface having a non-polar surface and a polar surface in which one of Si-polar surface and C-polar surface dominate. Another international patent application laid open No. WO2011/074237 has disclosed a technique to reduce defects formed at an interface between an SiC substrate and an insulating film by adding phosphorous (P) in the insulating film. Still another international patent application laid open No. WO2013/145023 has disclosed a technique to suppress variation in a threshold voltage of an FET by providing a laminated insulating film having a charge capturing function as the gate insulating film.\nA Japanese Patent Application laid open No. JP-2004-152813A has disclosed a device, which has a layer made of an SiC film of 3C-SiC type on an SiC substrate having a 6H-SiC type, which is operable in a high frequency by localizing electrons within the layer of 3C-SiC type. A Japanese Patent Application laid open No. JP-2013-197167A has disclosed a technique of forming a gate insulating film made of 3C-SiC type on an SiC substrate of 4H-SiC type. A conference report of Solid State Devices and Materials, Tokyo (1987) has disclosed a technique often called as the step controlling epitaxy where a lateral growth on a surface of an SiC substrate offset from a closest packed surface is carried out.\nA document of Nature Materials volume 8, pages 171 and 172 (2009) has disclosed a technique for growing a graphene layer on an SiC substrate by raising a temperature of the SiC substrate higher than 1100° C. by reducing the SiC substrate. An international patent application laid open No. WO2010/023934 has disclosed another technique for forming a graphene layer on an SiC substrate by exposing Si surface of the SiC substrate by removing a native oxide layer of the SiC substrate, oxidizing Si surface to form SiO2 layer, and heating the SiC substrate within a vacuum. A Japanese patent application laid open No. JP-2015-110485A and an international patent application laid open No. WO2013/125669 each have disclosed a technique to form a graphene layer by heating an SiC substrate within an inactive atmosphere to vaporize Si atoms. A Japanese patent application laid open No. JP-2013-510071A has disclosed a technique of hydrogenating a graphene layer.\nA Japanese patent application laid open No. JP-2014-162683A has disclosed a technique of removing hydrocarbons in a graphene layer by forming a carbon buffer layer on an SiC substrate by heating the SiC substrate, decomposing bonds between Si atoms in the SiC substrate and the carbon buffer by exposing the carbon buffer in hydrogen (H2) atmosphere, terminating Si bonds in the SiC substrate by hydrogen (H2), and heating the SiC substrate in a vacuum. A Japanese patent application laid open No. JP-2014-152051A has disclosed a technique of forming a graphene layer by removing native oxide on the surface of an SiC substrate to expose the C-polar surface of the SiC substrate, forming an SiC layer on the C-polar surface of the substrate, and heating the SiC substrate in argon (Ar) atmosphere. A Japanese patent application laid open No. JP-2013-180930A has disclosed a technique of forming a graphene layer by segregating carbon (C) atoms between an SiC substrate and an SiO2 film covering the SiC substrate to form a state where excess carbons are on the SiC substrate, and heating the SiC substrate at a temperature so as not to sublimate Si atoms. A Japanese patent application laid open No. JP-2014-240173A has disclosed a technique of forming a graphene layer by forming an SiC layer on a Si substrate, and heat-treating a surface of the SiC layer by using hydrogen (H2).\nThus, various trials and experiments have been carried out for enhancing the carrier mobility and reducing defects in an SiC-MOSFET. Also, various techniques have been proposed for forming a graphene layer having less defects and superior quality on an SiC substrate. However, no reports have seemed to be proposed how the SiC substrate has a preferable surface structure when a single crystal SiC substrate forms a film thereon."}
{"text": "As an example of a mobile communication system to which the present invention is applicable, a 3rd Generation Partnership Project Long Term Evolution (hereinafter, referred to as LTE) communication system is described in brief.\nFIG. 1 is a view schematically illustrating a network structure of an E-UMTS as an exemplary radio communication system. An Evolved Universal Mobile Telecommunications. System (E-UMTS) is an advanced version of a conventional Universal Mobile Telecommunications System (UMTS) and basic standardization thereof is currently underway in the 3GPP. E-UMTS may be generally referred to as a Long Term Evolution (LTE) system. For details of the technical specifications of the UMTS and E-UMTS, reference can be made to Release 7 and Release 8 of “3rd Generation Partnership Project; Technical Specification Group Radio Access Network”.\nReferring to FIG. 1, the E-UMTS includes a User Equipment (UE), eNode Bs (eNBs), and an Access Gateway (AG) which is located at an end of the network (E-UTRAN) and connected to an external network. The eNBs may simultaneously transmit multiple data streams for a broadcast service, a multicast service, and/or a unicast service.\nOne or more cells may exist per eNB. The cell is set to operate in one of bandwidths such as 1.25, 2.5, 5, 10, 15, and 20 MHz and provides a downlink (DL) or uplink (UL) transmission service to a plurality of UEs in the bandwidth. Different cells may be set to provide different bandwidths. The eNB controls data transmission or reception to and from a plurality of UEs. The eNB transmits DL scheduling information of DL data to a corresponding UE so as to inform the UE of a time/frequency domain in which the DL data is supposed to be transmitted, coding, a data size, and hybrid automatic repeat and request (HARM)-related information. In addition, the eNB transmits UL scheduling information of UL data to a corresponding UE so as to inform the UE of a time/frequency domain which may be used by the UE, coding, a data size, and HARQ-related information. An interface for transmitting user traffic or control traffic may be used between eNBs. A core network (CN) may include the AG and a network node or the like for user registration of UEs. The AG manages the mobility of a UE on a tracking area (TA) basis. One TA includes a plurality of cells.\nDevice to device (D2D) communication refers to the distributed communication technology that directly transfers traffic between adjacent nodes without using infrastructure such as a base station. In a D2D communication environment, each node such as a portable terminal discovers user equipment physically adjacent thereto and transmits traffic after setting communication session. In this way, since D2D communication may solve traffic overload by distributing traffic concentrated into the base station, the D2D communication may have received attention as the element technology of the next generation mobile communication technology after 4G. For this reason, the standard institute such as 3GPP or IEEE has proceeded to establish the D2D communication standard on the basis of LTE-A or Wi-Fi, and Qualcomm has developed their own D2D communication technology.\nIt is expected that the D2D communication contributes to increase throughput of a mobile communication system and create new communication services. Also, the D2D communication may support proximity based social network services or network game services. The problem of link of a user equipment located at a shade zone may be solved by using a D2D link as a relay. In this way, it is expected that the D2D technology will provide new services in various fields.\nThe D2D communication technologies such as infrared communication, ZigBee, radio frequency identification (RFID) and near field communications (NFC) based on the RFID have been already used. However, since these technologies support communication only of a specific object within a limited distance (about 1 m), it is difficult for the technologies to be regarded as the D2D communication technologies strictly.\nAlthough the D2D communication has been described as above, details of a method for transmitting data from a plurality of D2D user equipments with the same resource have not been suggested."}
{"text": "For a variety of reasons, a packet en route to its destination often passes through one or more networking device(s). For example, for security purposes, a packet may pass through various network devices, such as a firewall, a router, and an Intrusion Prevention System (IPS). Once a networking device has received the packet it may perform a series of complex operations on the packet before the packet is allowed to be transmitted or is dropped. For example, a firewall that has received a packet may perform a series of operations on the packet to determine whether the packet is valid and safe before transmitting it out the egress (output) interface.\nWith conventional approaches, a user may know that a packet has arrived in the ingress interface of a network device, but the user has no efficient way of determining the various operations performed on the packet within a network device. For example, with conventional approaches, a user may need to go through what could possibly be hundreds of logging messages to try to understand the various operations performed on a packet within a network device. In addition, in a typical network device, the conditions and rules that have been configured are generally numerous and complex, which means that the reason for dropping a packet in a network device may be elusive for a user on a time constraint and without full understanding of the device, or intimate knowledge of the protocols utilized.\nFor example, a traditional approach in analyzing the operations performed on the packet is by examining associated logs (syslogs) and debugs. However, a limitation of current logs and debugs is that the current logs and debugs do not apply to a single packet, but instead apply to an entire feature or protocol. As of result, trying to determine the operations performed on an individual packet by examining associated logs and/or debugs requires correlation of the logs/debugs back to the specific packet. This is very time consuming and inefficient."}
{"text": "1) Field of the Invention\nThe present invention relates to a technology for matching contents of a cache memory to contents of a database.\n2) Description of the Related Art\nSystems are known that make is possible to execute electronic-commerce transactions, over the Internet and using a web browser, from terminal of the users, or make it possible to access database of an enterprise, over the Internet and using a web browser, from terminal of the users. An application server is used in such systems as a middleware to control access to the database from the terminals of the users.\nAlthough it is difficult to handle status transition in the web browser, it is not so difficult in the application server. Therefore, the systems mentioned above can be easily construction and operates smoothly when the application server is provided as the middleware.\nIn such systems, a cache is provided to reduce frequency of access to the database. A related technology has been disclosed in, for example, Introduction to Web system, Jan. 23, 2003, HYPERLINK http://www.atmarkit.co.jp/fjava/rensai2/websys08/websys08.html.\nHowever, when the database is updated, a mismatch is produced between the contents in the cache memory and the contents in the database. To solve this problem, a notification that indicates that the database has been updated is sent to the application server and the application server refreshes the contents in the cache to match with the contents in the database.\nHowever, there is a problem that it is necessary to provide a means to generate and output the notification that indicates that the database has been updated.\nIn recent years, especially, parts-production of a program running in an application server or the like has advanced due to object-oriented languages such as JAVA 2 Enterprise Edition (J2EE). However, the parts-production of the program becomes difficult due to dependence on the processing performed at the database-side.\nFurther, cache contents in the application server are not stored at the database side, and a location in the database from which the cache contents have been extracted is not stored at the server side. Hence, when the database has been updated, an update notification is transmitted regardless of whether the cache contents should be refreshed, and moreover, the application server erases from the cache, even data that is not required to be updated. As a result, excessive cache refreshing occurs, which results in an increased load on the server.\nFurther, when the number of application servers connected to a database increases, load required for the database to perform a cache processing in each application server increases, which reduces the processing speed of the database."}
{"text": "As a measure against global warming, solar photovoltaic power generation has become popular in the world. For the solar photovoltaic power generation, a photoelectric conversion device (also called a solar cell) in which light energy is converted into electrical energy by using photoelectric characteristics of a semiconductor is applied in many cases, as compared to the case of utilizing solar heat.\nProduction of photoelectric conversion devices tends to increase year after year. For example, the total production of solar cells in the world in 2005 was 1,759 MW, which is a drastic increase of 147% as compared to that in the previous fiscal year. Photoelectric conversion devices which have become popular worldwide use a crystalline semiconductor; photoelectric conversion devices using a single crystal silicon substrate or a polycrystalline silicon substrate constitute a large part of the production.\nThe thickness of a crystal-type photoelectric conversion device using silicon, which is needed to absorb sunlight, is only about 10 μm. Nevertheless, a single crystal silicon substrate or a polycrystalline silicon substrate manufactured as a product has a thickness of about 200 to 300 μm. That is, the photoelectric conversion device using a single crystal semiconductor substrate or a polycrystalline semiconductor substrate has a thickness ten or more times as large as the thickness needed for photoelectric conversion, and thus the single crystal silicon substrate or the polycrystalline silicon substrate is far from being effectively utilized wholly. Speaking of extremes, most part of the single crystal silicon substrate or the polycrystalline silicon substrate functions only as a structural body for keeping the shape of the photoelectric conversion device.\nWith the increase of production of photoelectric conversion devices year after year, short of supply of polycrystalline silicon, which is the material of a silicon substrate, and resulting rise of cost of the same have become problems of the industry. The production of polycrystalline silicon is expected to be about 36,000 tons in 2007; in contrast, 25,000 tons or more of polycrystalline silicon is needed for semiconductor (LSI) and 20,000 tons or more of polycrystalline silicon is needed for solar cells, which means polycrystalline silicon seems to be in short of supply by about 10,000 tons. This short of supply is predicted to continue.\nThere are a variety of structures of photoelectric conversion devices. In addition to a photoelectric conversion device having a typical structure in which an n-type or a p-type diffusion layer is formed in a single crystal silicon substrate or a polycrystalline silicon substrate, a stacked photoelectric conversion device in which different kinds of unit cells, i.e., a unit cell formed of a single crystal semiconductor and a unit cell formed of an amorphous semiconductor, are combined is known (see Examined Patent Application Publication No. H6-044638). The photoelectric conversion devices are the same in that a single crystal silicon substrate or a polycrystalline silicon substrate is used. Here, as another mode of a photovoltaic device using a single crystal semiconductor substrate, a photovoltaic device using a single crystal semiconductor layer formed into a slice is given. For example, Patent Document 4 (Patent Document 4: Japanese Published Patent Application No. H10-335683) discloses a tandem solar cell in which hydrogen ions are implanted into a single crystal silicon substrate, a single crystal silicon layer which is separated from the single crystal silicon substrate in a layer shape is disposed over a support substrate in order to lower the cost and save resources while maintaining high conversion efficiency. In this tandem solar cell, a single crystal semiconductor layer and a substrate are bonded to each other with a conductive paste.\nOn the other hand, photoelectric conversion devices using a crystalline silicon thin film have also been developed. For example, a method for manufacturing a silicon thin-film solar cell in which a crystalline silicon film is deposited over a substrate by a plasma CVD method using a VHF of 27 MHz or higher which has been pulse modulated is described (see Japanese Published Patent Application No. 2005-50905). Further, a technique for controlling plasma process condition to optimize dopant concentration in crystal grains and crystal grain boundaries when a polycrystalline silicon thin film is formed by a plasma CVD method over a special electrode called a texture electrode which has minute unevenness on its surface is disclosed (see Japanese Published Patent Application No. 2004-14958). However, a crystalline thin-film silicon solar cell is still inferior to a single crystal silicon solar cell in crystal quality and photoelectric conversion characteristic. Moreover, a crystalline silicon film needs to be deposited to a thickness of 1 μm or more by a CVD method, which leads to a problem of low productivity."}
{"text": "Snappy is one algorithm in the LZ77 family that is widely used in applications such as Hadoop Compression, BTRFS file-system, in data center usages such as index-servers and is also enabled in the some browsers. In these usages, decompression latencies are critical to system performance."}
{"text": "1. Field\nEmbodiments relate to a ball balancer of a washing machine, and more particularly, to a relation between a gap, between a racer and balls, and viscous oil of a ball balancer of a washing machine.\n2. Description of the Related Art\nGenerally, washing machines wash laundry by rotating a drum containing the laundry using a driving motor. The drum is rotated at a low speed in regular and opposite directions during a washing process, and is rotated at a high speed in a regular direction during a dehydrating (drying) process.\nDuring the dehydrating (drying) process, the drum is rotated at a high speed under the condition that laundry is not equally disposed in all regions of the drum and is crowded at a specific region of the drum, or the laundry is pushed to one side of the drum due to the accelerated rotation of the drum at an initial state of the dehydrating (drying) process. As a result, the center of gravity of the drum does not coincide with the center of rotation of the drum, thus generating vibration and noise. When the above phenomenon is repeated, components of the washing machine, including a drum, a rotary shaft, and a driving motor, break down or have a shortened life span.\nParticularly, a drum washing machine has a structure in which a drum accommodating laundry is disposed horizontally so that the drum is rotated at a high speed in a dehydrating (drying) operation under the condition that the laundry is gathered together on the bottom of the drum by gravity. Thus, the center of gravity of the drum does not coincide with the center of rotation of the drum. Therefore, the drum washing machine has a great possibility of generating vibration and noise.\nDrum washing machines, in which a drum is disposed horizontally, and vertical axis washing machines, in which a drum is disposed vertically, are generally provided with balancers for maintaining the dynamic balance of the drum.\nKorean Patent Laid-open Publication No. 10-1999-0038279 discloses an example of a washing machine having balancers. Each of the balancers of this washing machine includes a racer installed at the upper or lower part of a drum for maintaining the dynamic balance of the drum when the drum is rotated at a high speed. Balls made of steel and freely movably are disposed in the racer, and viscous oil fills the inside of the racer.\nWhen the drum is rotated, the drum cannot maintain its dynamic balance due to the unbalanced eccentric structure of the drum and the partial distribution of laundry in the drum. Then, the steel balls compensate for the above unbalance, thus allowing the drum to maintain its dynamic balance.\nSince the viscous oil of the balancer employed by the above conventional washing machine is sensitive to an outdoor temperature, in the case that a gap between the inner wall of the racer and the steel balls is small, the ball balancer has a considerably high deviation in vibration (time) to attain a correct position according to the outdoor temperature.\nOn the other hand, in the case that the gap between the inner wall of the racer and the steel balls is large, the ball balancer rapidly attains a correct position, and thus the vibration of the drum is decreased. However, in this case, when the viscosity of the viscous oil is low, the ball balancer generates a large amount of noise. Further, when the gap is expanded to a certain extent, it is difficult to control vibration and noise characteristics due to the manufacturing error of the racer."}
{"text": "Various components used in microelectronic packaging or in interconnecting microelectronic packages can includes conductive features formed in a support structure in the form of a substrate or the like having vias or other routing or distribution circuitry formed therein. Such features are typically formed from various wiring metal such as copper, gold, nickel, aluminum, etc., which are formed in a support structure such as, for example, a thin dielectric or semiconductor material layer, by a number of different methods.\nIn some instances, the materials used to form such conductive features were made from a material of the type listed above in a form having impurities therein. Such materials, such as oxygen, nitrogen, chlorine, and carbon, occurred naturally within the material or were added to the material as a by-product of various mineral extraction or processing steps. These and other impurities within the wiring material were known to increase the resistivity of the wiring material, which can slow the performance of microelectronic packages or package assemblies including such wiring material and increase operating temperatures thereof. In addition, these impurities sometimes directly resulted in failure or the creation of defects within wiring circuitry by causing small explosions within the wiring material within which they were embedded. In an effort to reduce the resistivity of components and to otherwise increase their reliability, the previously discussed impurities were reduced in the wiring materials by additional processing, for example, to result in wiring material that was substantially more pure than that which was previously used. The removal of such impurities, however, lead to increased grain size within the wiring material used, which accordingly resulted in a decreased mechanical strength for such materials.\nThe reduced strength in the materials used to form wiring circuitry can lead to failure within the components within which they are formed. Such failure can include fracture within components due to heat cycling of support structures or other related components. Additional failure can be in the form of delamination from within support structures due to plastic deformation of the wiring components caused by elastic deformation of support structures or thermal expansion thereof. Accordingly, a wiring metal composition and associated method are needed that produce high strength wiring components while maintaining an acceptable level of resistivity without susceptibility to failure modes similar to those of low-purity wiring metals."}
{"text": "In the search for a suitable fuel to replace the high enriched UAlx fuel generally used in Research and Materials Test Reactors (MTR) with a lower enrichment fuel, one viable candidate is found to be a U—Mo alloy. Mo is added to metallic uranium to extend the stability domain of the high temperature gamma phase, since this phase is stable under the required irradiation conditions in contrast to the room temperature alpha phase. 7-10 wt % Mo is sufficient to avoid transformation to the alpha phase during the production process. U(Mo) particles are used in the fabrication of so-called dispersion fuel plates or rods, in which the U(Mo) particles are mixed with a matrix material (generally Al). In case of a fuel plate, the compacted powder mixture is subsequently pressed in between two Al alloy plates, after which this sandwich structure is rolled to the required thickness. In case of fuel rods, fabrication may often be based on coextrusion methods.\nA common method to manufacture the U(Mo) particles is by atomisation processes, which can best be described as a technique in which an ingot of U(Mo) alloy is molten using arc melting while it is spinning or a variation thereof. This causes molten material to be dispersed in small droplets by the centripetal forces. The droplets solidify on their way to the cooled chamber walls. The resulting spherical particles are very well suitable for the fabrication of the dispersion fuel plates or rods. Their sizes depend on the parameters used in the melting process, but are generally around 100 μm in diameter. Another production process based on grinding of ingots results in more or less irregular ground particles of similar average dimensions. Both production processes have been used in the production of test fuel plates.\nTests of U(Mo) based fuels in the reactor have revealed that the U(Mo) particles (atomised and ground) form an interaction layer with the Al matrix under irradiation. This interaction layer was demonstrated to amorphise under irradiation and shows very poor fission gas retention. This causes swelling of the fuel plates, which eventually leads to their failure by pillowing. In an attempt to remedy this behaviour, it was found that Si addition to the Al matrix results in some improvement, although it may not provide a complete solution, particularly for higher power densities. In “Dispersed (Coated particles) and monolithic (zircalloy-4 cladding) U—Mo Particles” (RERTR—2005 Meeting), Pasqualini describes CVD coating of U(Mo) particles with a silicon coating for introducing silicon as an inhibitor. Chemical vapour deposition of silicon thereby is based on silane. Furthermore, silicon is undesireable for the reprocessing of the used fuel and an alternative way, without Si, to stabilise the behaviour of the U(Mo) fuel during irradiation would be beneficial.\nThe incorporation of neutron poisons in nuclear fuels is a frequently used method to fine-tune the characteristics of the fuel towards its use in the reactor. Neutron poisons typically may have the purpose of lowering the reactivity at the beginning of the use of a fuel assembly in the reactor. Incorporation of neutron poisons is preferably homogeneous throughout the fuel and is often accomplished by blending neutron poison powders in the matrix. Because of several reasons, eg. the higher uranium loadings required in the LEU based fuels, this blending method is less appropriate for LEU based fuels and efforts are made to incorporate the neutron poison in the structural materials of the fuel element. In “Cd wires as burnable poison for the BR2 fuel element” by Franck et al. in Transactions of RRFM 2009, it is suggested to replace the provision of neutron poisons as powder mixed in the matrix by the provision of wires of neutron poison material in the structural materials of the fuel element."}
{"text": "When professionals develop training applications using eLearning tools, they typically create such applications section-by-section. The creation of each section often includes setting background and other art and also implementing user-interactive attributes (e.g., buttons that represent choices and cause an application to jump to different sections) in some of the sections. If later during the development process, the professional desires to go back and change one or more interactions in a previous section, the professional may have to search section-by-section until he or she finds the particular interaction, then edit the interaction in the section creation interface. This can be very burdensome, especially for large applications with many sections because using a “poke-and-peek” strategy may require looking through an unreasonable number of previous sections before the interaction is found.\nOne prior art solution has been to manually create a story board during the development process, wherein each section is laid out with its characteristics and interactions listed. A disadvantage of this solution is that it uses time and effort for something that is purely administrative and takes time and effort away from the development process. There is no solution currently available that allows a user to view and access interactions from a number of sections in a convenient and well-organized manner, thereby dispensing with the current section-by-section search approach."}
{"text": "This invention relates to methods and apparatus for deriving substantially correct parameter data for processing workpieces, more particularly, processing workpieces for electronic device fabrication.\nThe most successful processing of materials for electronic devices typically requires optimization and precise control of the processing environment at all process steps. Many of these process steps are performed under conditions that make it difficult or impossible to measure the desired process variables. In those cases where an important process variable cannot be readily measured, an attempt is made to correlate the parameter of interest to other measurable or controllable parameters. The accuracy and stability of these correlations, also called equipment response models, are a critical factor in determining the process capability and device yield at any given process step.\nDescriptions of some of the available technologies and sensor apparatus for measuring process variables are available in the technical and patent literature. Examples of some of the technologies are described in United States patents U.S. Pat. Nos. 6,244,121, 6,051,443, 6,033,922, 5,989,349, 5,967,661, 5,907,820, and Patent Corporation Treaty application WO 02/17030.\nSome of the available technologies are tethered systems in which sensors exposed to the process conditions that are to be measured have physical connections to remote facilities such as power sources and information processors, and electronic components. Other technologies use an electronics module that is coupled to the detectors on a support; the sensors and electronics module are part of a single unit that is exposed to the process conditions that are to be measured. The electronics module is necessary for a sensor apparatus with capabilities such as autonomous information processing capability, wireless communication capability, and other electronically controlled on-board capabilities.\nIn some applications, the presence of the tether or electronics module can introduce an unacceptable error in the measurement data. For the case of measuring temperatures using a sensor apparatus, the mere presence of the electronics module will distort the temperature field being measured. It is possible to reduce the distortion by using very small components in the module, thereby reducing the module\"\"s overall size and thermal mass.\nOf course, for some applications the module distortion effect may be ignored if a high degree of measurement accuracy is unnecessary. However, some of the critical process steps required for processing high-value workpieces such as semiconductor wafers for electronic device and optical device fabrication and substrates for flatpanel display fabrication do indeed require high accuracy for the values of the process conditions. For such applications, measurements of a parameter such as temperature need to be extremely accurate, particularly for the temperature uniformity across the area of the workpiece. Furthermore, applications involved with the manufacture, calibration, research, and development of process equipment for processing high value substrates require high accuracy measurements since the operation of the equipment can be limited by the accuracy of the measurements. Inaccurate data can result in the loss of millions of dollars of product in some instances. Alternatively, the inaccurate data can result in the production of products having poorer performance because the process conditions were optimized based on the inaccurate data.\nClearly, there are numerous applications requiring high accuracy methods and apparatus by which spatially resolved and time resolved equipment response models can be easily and economically developed and maintained. An example of an important application is the uniform processing of workpieces such as semiconductor wafers, flatpanel displays, and other electronic devices. Furthermore, there is a need for high accuracy methods and apparatus capable of collecting data for response models in a nonperturbing manner on unmodified process equipment running realistic process conditions. Still further, there is a need for methods and apparatus capable of correcting measurements errors that can be caused by the methods and apparatus used for the measurements.\nThis invention seeks to provide methods and apparatus that can improve the accuracy of measured parameter data used for processing workpieces. One aspect of the present invention includes methods of deriving substantially correct data for applications such as generating data for monitoring, controlling, and optimizing processes and process tools. Another aspect of the present invention includes apparatus for deriving substantially correct data for applications such as generating data for monitoring, controlling, and optimizing processes and process tools.\nIt is to be understood that the invention is not limited in its application to the details of construction and to the arrangements of the components set forth in the following description or illustrated in the drawings. The invention is capable of other embodiments and of being practiced and carried out in various ways. In addition, it is to be understood that the phraseology and terminology employed herein are for the purpose of description and should not be regarded as limiting.\nAs such, those skilled in the art will appreciate that the conception, upon which this disclosure is based, may readily be utilized as a basis for the designing of other structures, methods and systems for carrying out aspects of the present invention. It is important, therefore, that the claims be regarded as including such equivalent constructions insofar as they do not depart from the spirit and scope of the present invention.\nThe above and still further features and advantages of the present invention will become apparent upon consideration of the following detailed descriptions of specific embodiments thereof, especially when taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display (LCD) device, and more particularly, to an LCD device that performs a gradation display by a combination of a plurality of sub frames according to a gradation level represented by a plurality of bits.\n2. Description of the Related Art\nIn the past, a sub frame drive scheme has been known as one of half-tone display schemes used in an LCD device. In the sub frame drive scheme which is a kind of a time axis modulation scheme, a certain period of time (for example, a frame which is a display unit of an image in a case of a moving image) is divided into a plurality of sub frames, and the sub frames are combined according to a gradation to be displayed so as to drive pixels. The gradation to be displayed is decided according to a ratio of a pixel drive period of time occupied in the certain period of time. The ratio of the pixel drive period of time occupied in the certain period of time is decided according to a combination of the divided sub frames.\nAs an LCD device employing the sub frame drive scheme, known is an LCD device including pixels each of which includes a master latch, a slave latch, an LCD element, and three switching transistors of first to third switching transistors as disclosed in, for example, Japanese Translation of PCT International Application Publication No. JP-T-2001-523847 is known. In this case, in each pixel, first data of one bit is applied to one of two input terminals of the master latch through the first switching transistor; and second data of one bit having a complementary relation with the first data is applied to another input terminal through the second switching transistor. Then, when a target pixel is selected as a row selection signal applied through a row scanning line, the first switching transistor and the second switching transistor are turned on, and the first data is written. When the first data has a logical value of “1” and the second data has a logical value of “0,” the pixel performs a display based on data.\nAfter each data is written in all pixels by the above-described operation within a certain sub frame period of time, the third switching transistors of all pixels are turned on within the sub frame period of time. Then, the data written in the master latch is read out to the slave latch with a certain time difference. Then, the data latched in the slave latch is applied to a pixel electrode of the LCD element. The above-described series of operations are repeated at intervals of sub frames, and a desired gradation display is performed based on combinations of all sub frames within one frame period of time.\nIn other words, in the LCD device employing the sub frame drive scheme, certain display periods of time which are the same as or different from each other are allocated to all sub frames present within one frame period of time. Further, each pixel performs a white display in all sub frames at the time of a maximum gradation display (displayed), but does not perform a white display in all sub frames at the time of a minimum gradation display (not displayed, that is, displayed black). In the case other than at the time of the maximum gradation display or the minimum gradation display, a sub frame to be displayed white is selected according to a gradation to be displayed white. In addition, in the LCD device according to the related art, data to be input is digital data representing a gradation, and a digital drive scheme of a two-stage latch configuration is used.\nHowever, in the LCD device according to the related art, each of the two latches in each pixel is configured with a so-called static random access memory (SRAM), and thus the number of transistors configuring the circuit becomes large. Thus, there is a problem in that it is difficult to reduce the size of a pixel. Further, data written in the two-stage master latches is read out to the slave latch at the same time, and the data latched in the slave latch is applied from the slave latch to the pixel electrode of the LCD element. At this time, however, when switching is performed on all pixels at the same time (at one time), all pixels are read out and thus a consumption current at the moment becomes extremely high. The occurrence of an instantaneous peak of consumption current results in a drop in a power voltage or an increase in a GND voltage, which seriously influences the drive operation of the whole LCD device."}
{"text": "(a) Field of the Invention\nThe present invention is concerned with a process and a device for determining the quantity of impurities in a gas by chromatography in gas phase and the utilization of such a process for measuring the quantity of doping impurities found in traces in silane.\n(b) Description of Prior Art\nThe technology of chromatography, under normal conditions, does not permit detection to reach thresholds lower than about 2.10.sup.-2 ppm.volume and it may be necessary to determine the quantity of impurities which are found in a gas as traces, in lower quantities. This is particularly the case during the utilisation of silane in the manufacture of semi-conductors. This gas should indeed contain so called doping gas impurities, such as phosphine (PH.sub.3), arsine (AsH.sub.3) or hydrogen sulfide (H.sub.2 S). The effect of these impurities once incorporated into the semi-conductor should be controlled and maintained at a very specific level of concentration, which is very low, is exceeded. For example, it may be necessary to detect in silane traces of phosphine and arsine which are not in excess of 10.sup.-5 ppm.volume in order to be able to guarantee the quality of the silane used in the manufacture of semi-conductors."}
{"text": "This disclosure relates to a ring gear used in an epicyclic gear train of a gas turbine engine.\nGas turbine engines typically employ an epicyclic gear train connected to the turbine section of the engine, which is used to drive the turbo fan. In a typical epicyclic gear train, a sun gear receives rotational input from a turbine shaft through a compressor shaft. A carrier supports intermediate gears that surround and mesh with the sun gear. A ring gear surrounds and meshes with the intermediate gears. In arrangements in which the carrier is fixed against rotation, the intermediate gears are referred to as “star” gears and the ring gear is coupled to an output shaft that supports the turbo fan.\nTypically, the ring gear is connected to the turbo fan shaft using a spline ring. The spline ring is secured to a flange of the turbo fan shaft using circumferentially arranged bolts. The spline ring includes splines opposite the flange that supports a splined outer circumferential surface of the ring gear. The ring gear typically includes first and second portions that provide teeth facing in opposite directions, which mesh with complimentary oppositely facing teeth of the star gears.\nAn epicyclic gear train must share the load between the gears within the system. As a result, the splined connection between the ring gear and spline ring is subject to wear under high loads and deflection. Since the spline connection requires radial clearance, it is difficult to get a repeatable balance of the turbo fan assembly. Balance can also deteriorate over time with spline wear."}
{"text": "1. Field of the Invention\nThis invention relates to laser technology, and more particularly, to a polarization-type laser detection system which is capable of distinguishing between a reflected light signal from a target object and a back-scattered light signal from a mass of suspended particles in the air.\n2. Description of Related Art\nFIG. 1 is a schematic diagram of a conventional laser detection system. As shown, this laser detection system includes a control unit 10, a laser emitter 20, and an optical receiver 30. This laser detection system is used to detect whether any target object is present nearby. In the case of a target object 40 is present nearby, the laser beam from the laser emitter 20 will be reflected back by the target object 40 and received by the optical receiver 30. The received light is then analyzed by the control unit 10 to indicate the presence of the target object 40.\nOne drawback to the foregoing laser detection system, however, is that, under bad weather conditions when the air is filled with suspended particles, the emitted laser beam from the laser emitter 20 would be scattered back, causing the optical receiver 30 to received a back-scattered light signal that would make the control unit 10 unable to perform the intended target detection. The laser detection system would therefore operate improperly under bad weather conditions."}
{"text": "1. Field of the Invention\nThe present invention relates to a plasma lamp system, and more particularly, to a plasma lamp system and a bulb therefor capable of maximizing lighting efficiency by improving a point source of light characteristic and a spectrum characteristic.\n2. Description of the Background Art\nIn general, a plasma lamp system is a lighting system which emits visible rays and ultraviolet rays when a filling material within a bulb is excited by microwave energy or electric discharge, has a long life span compared to an incandescent lamp or a fluorescent lamp, and has excellent efficiency in lighting.\nA bulb of a plasma lamp system is filled with high-pressure mercury or metal halide as a primary light-emitting material leading light-emitting when excited by microwave energy or electric discharge, together with inert gas for forming a plasma at an initial stage of lighting-emitting, additives which make lighting easy and the like.\nSuch a conventional plasma lamp system shows different characteristics according to types of the primary light-emitting materials within the bulb.\nIn a UHP (ultra high performance) lamp system having a bulb filled with high-pressure mercury, mercury of about 200 atm or more emits light in operation and shows a spectrum characteristic as shown in FIG. 1. That is, in the UHP lamp system, intensity of light is high in a blue color region where a wavelength of light is about 400˜500 nanometers and in a green color region where a wavelength of light is about 500˜600 nanometers. But, intensity of light is low in a red color region where a wavelength of light is about 600˜700 nanometers. Accordingly, the UHP lamp system cannot obtain high reddish color purity, thereby having limitations on improving display performance.\nIn an MH (metal halide) lamp system having a bulb filled with metal halide, since intensity of light shows an independent peak in each of red, green and blue color regions, it is easy to make a proper optical spectrum. But, optimum pressure of the metal halide is relatively low because of a characteristic of metal halide, and thus the light is emitted not in parallel but radially, whereby it is not easy to apply the MH lamp system to an optical system that requires a point source of light. Therefore, in case of applying the MH lamp system to an optical system which requires a point source of light, such as a projector or a projection display, as shown in FIG. 2, the center of an arc 40 of a bulb 10 should be positioned at a focal point of a reflector 20 for reflecting light in order to obtain parallel light 30, thereby causing intricacy in manufacturing and degrading productivity."}
{"text": "1. Technical Field\nThe present invention relates to pattern recognition and object recognition techniques and, more particularly, to a human face detection position deviation-correcting method, correcting system, and correction program capable of detecting whether a detected human face has deviated in position and of correcting the position.\n2. Related Art\nIn recent years, as pattern recognition technology and information-processing apparatuses such as computers have improved in performance, the accuracy at which characters and voice are recognized has improved at an amazing rate. It is known, however, that it is still difficult to recognize persons, objects, and scenes within images taken, for example, by a still image camera by pattern recognition technology, and that it is especially quite difficult to precisely and quickly recognize whether a human face is contained in such an image.\nHowever, it has become quite important to recognize whether a human face is contained in such an image and to identify who is in the picture precisely and automatically using a computer in establishing biometric recognition technology, improving security, speeding up criminal investigations, increasing the rates at which image data are put in order, and speeding up search operations.\nMeanwhile, in the case of ordinary users who are fond of photography and designers professionally engaging in image processing, they frequently want to make layouts about the subject human face from the taken digital photograph. In one case, for example, a picture where the face has been zoomed in is created from a picture containing the entire human body. In another case, a picture for a personal history is created.\nIn such cases, it is necessary to zoom into the face while precisely aligning the center of the face or to align the center of the face to the center of the photograph frame manually. Hence, complex and skilled operations are necessary.\nTherefore, in the past, a method of detecting the position of a face using a facial image detection technique and aligning the position of the facial image has been proposed.\nFor example, in H. A. Rowley, S. Baluja, and T. Kanade, “Neural network-based face detection”, IEEE Transactions on Pattern Analysis and Machine Intelligence, Vol. 20, No. 1, pp. 23-38, 1998, there is disclosed a technique for detecting a facial image by training a neural network using facial images.\nFurthermore, in the above paper and so on, a decision is first made as to whether there is any region of human flesh color within an input image. A mosaic size is automatically determined about the region of human flesh color. A mosaic candidate region is created. The proximity to the dictionary of human faces is calculated. Thus, a decision is made as to whether there is a human face. The face is extracted. This reduces erroneous extraction due to the effects of the background. Consequently, human faces are found from images automatically and efficiently.\nWith the method of the above-cited paper, it is possible to detect a facial image but there is the problem that the center of the position of a human face can not always be precisely calculated as disclosed in the paper.\nOn the other hand, in JP-A-H9-171560, it is necessary to calculate the center line of the rotation while creating data about rotation of the face. Hence, there is the anxiety that the processing time is increased. Furthermore, it is assumed that every assumed pixel has symmetry. Therefore, it is quite difficult to stably gain the center line due to noise.\nAccordingly, the present invention has been devised to effectively solve these problems. It is an object of the present invention to provide a novel human face detection position deviation-correcting method, correcting system, and correcting program capable of judging whether a detected human face has deviated from the detection position and of correcting the positional deviation from the detection position by a quick and simple method."}
{"text": "1. Field of the Invention\nThe present invention relates in general to the field of plastics. More particularly, the present invention relates to a pallet system and a pallet apparatus and a method of pallet manufacturing. Specifically, one preferred embodiment of the present invention relates to blow molding a first pallet piece and then inserting it into another mold to join it to a second blow molded pallet piece.\n2. Discussion of the Related Art\nPallets are used for transporting and storing various types of cargo throughout the world. Therefore, it is important for a pallet to be lightweight, compact and strong. Traditionally pallets were made of wood, however, plastic pallets are now being used with increased frequency as they are lightweight, generally stronger, do not absorb moisture, and they can be cleaned and disinfected which is important in particular for the food industry. One disadvantage of plastic pallets is that they tend to break up with continued use over time. Therefore, for these and other various reasons it is important that plastic pallets be reinforced by a variety of means.\nThe below-referenced U.S. patents disclose embodiments that were at least in-part satisfactory for the purposes for which they were intended. The disclosures of all the below-referenced prior United States patents, in their entireties, are hereby expressly incorporated by reference into the present application for purposes including, but not limited to, indicating the background of the present invention and illustrating the state of the art.\nU.S. Pat. No. 5,417,167 discloses a plastic shipping pallet of the type usually manufactured from wood which is made of hollow plastic stringers and deck boards. The stringers and deck boards are made of plastic materials that may be either virgin plastic, recycled plastics or mixes. The stringers and deck boards are hollow but have closed ends to prevent entry of dirt, fluids, insects and vermin. The stringers and deck boards may be made in a blow molding process or using continuous extrusion and molding processes similar to that used in the manufacture of corrugated pipe. Various connection configurations for affixing the deck boards to the stringers are described including fixation by screws as well as interlocking connections between the plastic parts themselves and the use of separate interlocking components which engage the deck board and stringers.\nU.S. Pat. No. 5,845,588 discloses a pallet fabricated of a thermoplastic material which is produced by joining together individually fabricated structural components, each of which is individually thermoformed from an extruded plastic parison of a multilayer structure. The structural components include a deck in the form of a continuous sheet with longitudinal and transverse sides, one or more upper runner components fastened by their top side to the bottom of the deck, each of which includes a cross member with downward-facing posts, and one or more lower runner components, each of which includes a cross member with upward-facing posts and which is fastened by the tops of its posts to the bottoms of matching posts on the corresponding upper runner components.\nU.S. Pat. No. 5,868,080 discloses a reinforced plastic pallet construction and assembly method wherein multiple reinforcing bars are employed. At least some of the reinforcing bars have an exposed surface at a topside surface, underneath surface or underside surface of the pallet. In addition to functioning as a reinforcing member, the exposed surfaces of the reinforcing bars comprise an anti-skid surface for maintaining positioning of payload on the pallet or facilitating transport of the pallet, e.g., via a forklift or automated transport system. Various techniques for retaining the reinforcing bars within channels formed in the plastic pallet body are described. The reinforcing bars preferably comprise composite structural members of fiberglass reinforced plastic fabricated from a pultrusion process.\nU.S. Pat. No. 6,209,464 discloses a pallet that includes a rectangular support deck having a substantially planar upper supporting surface including a plurality of channels formed therein and extending across the pallet. The deck includes support members extending from an underside of the deck that nest in recesses formed in the upper surface of the deck. An alignment portion is formed at a center point along the first edge of the deck and a second alignment portion is formed at a center point of an opposed edge of the deck. The alignment portion includes angled sides for receiving a tapered member of a complementary device."}
{"text": "The Service Availability™ (SA) Forum has developed standard interfaces to enable the delivery of highly available carrier-grade systems with off-the-shelf platforms, middleware and service applications. Implementation of the standards allows software developers to focus on the application code that provides mission-critical services, and to minimize the need for customized high availability programming.\nFIG. 1 is a conceptual architecture stack defined by the SA forum for a highly available (HA) system. The SA forum has developed the Application Interface Specification (AIS) 120 to provide a standardized interface between HA applications 110 and HA middleware 130, such that they can be independent of one another. The HA applications 110 and the HA middleware 130 are run on an operating system 140 and are hosted by a hardware platform 150. Service continuity is achieved with the cooperation of all of the components in the stack. As part of the AIS 120, the SA forum has also developed a specification for the Availability Management Framework (AMF) 125 to coordinate and manage service availability in an HA system (SA Forum, Application Interface Specification, Availability Management Framework SAI-AIS-AMF-B.04.01). An AMF configuration, which is a configuration compliant with the AMF specification, is a logical organization of resources for providing and protecting services for a given HA application. An AMF configuration consists of components grouped into service units (SUs), which are further grouped into services groups (SGs). An application may consist of one or several SGs to provide and protect services. The services are defined in terms of service instances (SIs) composed of component service instances (CSIs). At runtime, for each SI, the middleware implementing the AMF service assigns active and standby HA states to SUs, according to a redundancy model used by the HA system. The AMF specification defines five redundancy models, which include the No-Redundancy, 2N, N+M, N-Way-Active and N-Way redundancy models.\nSeveral approaches are devised for generating valid AMF configurations. The design of the AMF configuration is driven by the required service and the specified redundancy model. Some of the approaches systematically generate AMF configurations that ensure SI protection: for example, components are created to provide and protect requested CSIs, and SUs are defined to provide and protect required SIs. In addition, the SUs are grouped into SGs with the redundancy model requested by the designer. However, when configurations are designed in an ad hoc manner, there can be a question of whether the configurations ensure SI protection. In some situations, new or additional services may be requested after a configuration is generated. In these situations, there can be a question of whether the current configuration provides and protects all the services (including the new/additional services).\nAMF configurations created by designers using traditional tools also need to be validated. Validating the protection of services as requested may require the exploration of all possible SI combinations and SI-SU assignments. This is a combinatorial and complex problem. In general, the problem is NP-hard for the N+M, N-Way-Active, and N-Way redundancy models."}
{"text": "Many ribonucleic acid (RNA) molecules contain secondary structure that results from hybridization between complementary regions within the RNA molecule. A variety of secondary structures can be formed, including hairpins and cruciforms. RNA molecules containing secondary structure are often difficult to reverse transcribe because polymerases cannot readily process through the secondary structure.\nBecause of the difficulty of reverse transcribing RNA molecules with secondary structure, many techniques dependent upon reverse transcription yield anomalous results. For example, RNA molecules with secondary structure may be poorly represented in cDNA libraries. Populations of RNA with secondary structure may also yield cDNA libraries with a short insert size. Furthermore, RNA molecules containing secondary structure may be difficult to detect in assays such as reverse transcription-polymerase chain reaction (RT-PCR).\nTraditionally, reverse transcription has been performed with reverse transcriptases encoded by retroviruses (e.g., avian myoblastosis virus (AMV) reverse transcriptase and Moloney murine leukemia virus (MMLV) reverse transcriptase). Several mesophillic DNA polymerases (e.g., E. coli DNA polymerase I) have also been shown to possess reverse transcriptase activity. However, these enzymes are generally used at temperatures of between about 37° C. to 42° C., a temperature range where secondary structure can be a significant problem.\nSeveral thermophilic DNA polymerases (e.g., Thermus aquaticus DNA polymerase and Thermus thermophilus DNA polymerase) also have reverse transcriptase activity. These enzymes are useful for reverse transcription, because at the high temperatures where such enzymes are stable, secondary structure in RNA molecules is reduced. Furthermore, such enzymes can be used to directly synthesize second strand DNA and potentially even to directly amplify an RNA target. However, the utility of these thermostable enzymes is limited because they require manganese as a co-factor for reverse transcriptase activity (e.g., U.S. Pat. No. 5,322,770) resulting in deleterious effects. In some cases, the fidelity of the polymerase is reduced as compared to the fidelity of the enzyme in the presence of other cofactors, such as magnesium ions. Therefore, it is not desirable to amplify the template in the same reaction mixture in which reverse transcription reaction is conducted. This necessitates extra time consuming steps when performing RT-PCR. In other cases, the presence of manganese ions may also cause degradation of the RNA template.\nAccordingly, what is needed in the art are alternative thermostable polymerases that have reverse transcriptase activity. Preferably, such thermostable polymerases should have reverse transcriptase activity in the presence of magnesium so that high-fidelity cDNAs may be obtained and so that both reverse transcription and amplification in RT-PCR reactions may conducted in the same reaction mixture."}
{"text": "Soybean (Glycine max), is an important and valuable field crop and is the world's leading source of vegetable oil and protein meal. Thus, there is a continuing need to develop new stable, high yielding soybean cultivars that are agronomically sound. New soybeans cultivars are desired that exhibit modified phenotypical traits, including modification of growth, temperature requirements, initiation date of floral or reproductive development, increased seed yield, resistance to diseases, nematodes and insects, abiotic stress tolerance, production of better stems and roots, altered fatty acid profiles, improvements in compositional traits, and better agronomic quality. One approach to produce such new soybean cultivars is through traditional breeding techniques wherein the soybean breeder selects and develops soybean plants having traits that result in superior varieties.\nThere are numerous steps in the development of any novel, desirable plant germplasm. These processes, which lead to the final step of marketing and distribution, can take from six to twelve years from the time the first cross is made. Therefore, development of new cultivars is a time-consuming process that requires precise forward planning and efficient use of resources. Furthermore, due to the number of genes within each chromosome, millions of genetic combinations exist in the breeders' experimental soybean material. Thus, developing a single variety of useful commercial soybean germplasm is highly unpredictable, and requires intensive research and development. The present invention relates to a new and distinctive soybean cultivar, designated TMG4177."}
{"text": "The invention relates generally to the fabrication of semiconductor devices and, more particularly, to the fabrication of high-K dielectric layers in semiconductor devices.\nFabrication of semiconductor devices, such as a metal-oxide-semiconductor (MOS) integrated circuit, involves numerous processing steps. In a semiconductor device, a gate dielectric, typically formed from silicon dioxide (xe2x80x9coxidexe2x80x9d), is formed on a semiconductor substrate which is doped with either n-type or p-type impurities. For each MOS field effect transistor (MOSFET) being formed, a gate conductor is formed over the gate dielectric, and dopant impurities are introduced into the substrate to form a source and drain. A pervasive trend in modern integrated circuit manufacture is to produce transistors having feature sizes as small as possible. Many present processes employ features, such as gate conductors and interconnects, which have less than 0.18 xcexcm critical dimension. As feature sizes continue to decrease, the size of the resulting transistor as well as the interconnect between transistors also decreases. Fabrication of smaller transistors allows more transistors to be placed on a single monolithic substrate, thereby allowing relatively large circuit systems to be incorporated on a single, relatively small die area.\nAs MOSFET feature sizes decrease, gate oxide thickness decreases as well. This decrease in gate oxide thickness is driven in part by the demands of overall device scaling. As gate conductor widths decrease, for example, other device dimensions must also decrease in order to maintain proper device operation. Early MOSFET scaling techniques involved decreasing all dimensions and voltages by a constant scaling factor, to maintain constant electric fields in the device as the feature size decreased. This approach has given way to more flexible scaling guidelines which account for operating characteristics of short-channel devices. For example, a maximum value of MOSFET sub-threshold current can be maintained while feature sizes shrink, by decreasing any or all of several quantities, including gate oxide thickness, operating voltage, depletion width, and junction depth, by appropriate amounts.\nAs a result of the continuing decrease in feature size, gate oxide thickness has been reduced so much that oxides are approaching thicknesses on the order of ten angstroms (xc3x85). Unfortunately, thin oxide films may break down when subjected to an electric field, particularly for gate oxides less than 50 xc3x85 thick. It is probable that even for a relatively low gate voltage of 3V, electrons can pass through such a thin gate oxide by a quantum mechanical tunneling effect. In this manner, a tunneling current may undesirably form between the semiconductor substrate and the gate conductor, adversely affecting the operability of the device. It is postulated that some of these electrons may become entrapped within the gate oxide by, e.g., dangling bonds. As a result, a net negative charge density may form in the gate oxide. As the trapped charge accumulates with time, the threshold voltage VT may shift from its design specification. Breakdown of the gate oxide may also occur at even lower values of gate voltage, as a result of defects in the gate oxide. Such defects are unfortunately prevalent in relatively thin gate oxides. For example, a thin gate oxide often contains pinholes and/or localized voids due to unevenness at which the oxide grows on a less than perfect silicon lattice.\nA more promising approach to further increasing gate dielectric capacitance may be to increase the permittivity of the gate dielectric. Permittivity, xcex5, of a material reflects the ability of the material to be polarized by an electric field. The permittivity of a material is typically described as its permittivity normalized to the permittivity of a vacuum, xcex50. Hence, the relative permittivity, referred to as the dielectric constant, of a material is defined as:\nK=xcex5/xcex50 \nWhile silicon dioxide (sometimes simply referred to as xe2x80x9coxidexe2x80x9d) has a dielectric constant of approximately 3.9, other materials have higher K values. Silicon nitride (xe2x80x9cnitridexe2x80x9d), for example, has a K of about 6 to 9 (depending on formation conditions). Much higher K values of, for example, 20 or more can be obtained with various transition metal oxides including tantalum oxide (Ta2O5), barium strontium titanate (xe2x80x9cBSTxe2x80x9d), and lead zirconate titanate (xe2x80x9cPZTxe2x80x9d). Using a high-K dielectric material for a gate dielectric would allow a high capacitance to be achieved even with a relatively thick dielectric layer. For example, a nitride gate dielectric with a K of 7.8 and a thickness of 100 angstroms is substantially electrically equivalent to an oxide gate dielectric (K about 3.9) having a thickness of about 50 angstroms. For even higher-K dielectric materials, even thicker gate dielectric layers could be formed while maintaining capacitance values higher than are possible with even very thin oxide layers. In this way, the reliability problems associated with very thin dielectric layers may be avoided while transistor performance is increased.\nOne problem which has been reported relating to integration of high-K dielectric materials is oxidation of silicon by certain high-K dielectric materials when the high-K dielectric material is formed directly on a silicon substrate. Since oxidation results in formation of what may be referred to as a xe2x80x9cstandard-Kxe2x80x9d dielectric material, some of the benefit of the high-K dielectric material is considered to be lost. In addition, reactions considered adverse between the high-K dielectric material and standard-K dielectric materials may also occur.\nThus, a method of forming a relatively high-K dielectric material which either overcomes or takes advantage of such reactions, and which provides the electrical advantages of a higher K, is needed.\nThe present invention relates to a method of making a semiconductor device having a composite dielectric layer, comprising: providing a semiconductor substrate; depositing on the semiconductor substrate alternating sub-layers of a first dielectric material and a second dielectric material to form a layered dielectric structure having at least two sub-layers of the first dielectric material and at least one sub-layer of the second dielectric material, wherein the first dielectric material is a high-K dielectric material and the second dielectric material is a standard-K dielectric material, and at least one of the one or more dielectric material sub-layers contain nitrogen implanted therein using a nitridation step; and annealing the layered dielectric structure at an elevated temperature to form a composite dielectric layer about the boundary of each first dielectric material layer/second dielectric material layer.\nThe present invention further relates to a method of making a semiconductor device having a composite dielectric layer, comprising: providing a semiconductor substrate; depositing on the semiconductor substrate alternating sub-layers of a first dielectric material and a second dielectric material to form a layered dielectric structure having at least two sub-layers of the first dielectric material and at two sub-layers of the second dielectric material, wherein the first dielectric material is a standard-K dielectric material and the second dielectric material is a high-K dielectric material, and at least one of dielectric material sub-layers contain nitrogen implanted therein using a nitridation step; and annealing the layered dielectric structure at an elevated temperature to form a composite dielectric layer about the boundary of each first dielectric material layer/second dielectric material layer.\nThe present invention further relates to a method of making a semiconductor device having a composite dielectric layer, comprising: providing a semiconductor substrate; subjecting the semiconductor substrate to a nitridation step to produce a layer of standard-K dielectric material in the upper portion of one side of the semiconductor substrate; depositing on the standard-K dielectric side of the semiconductor substrate alternating sub-layers of a first dielectric material and a second dielectric material to form a layered dielectric structure having at least one sub-layer of the first dielectric material and at least one sub-layer of the second dielectric material, wherein the first dielectric material is a high-K dielectric material and the second dielectric material is a standard-K dielectric material, and at least one of the one or more dielectric material sub-layers contain nitrogen implanted therein using a nitridation step; and annealing the layered dielectric structure at an elevated temperature to form a composite dielectric layer about the boundary of each first dielectric material layer/second dielectric material layer.\nThus, the present invention overcomes the problem of forming a high-K dielectric material which overcomes and takes advantage of previously disfavored reactions between dielectric materials, to form a composite dielectric layer which includes a reaction product of the high-K dielectric material and the standard-K dielectric material, which is obtained by annealing a layered dielectric structure at an elevated temperature."}
{"text": "Integrated circuit devices generally include an integrated circuit chip and a lead frame which are sealed within a protective enclosure, find wide use in various products, among which are consumer electronics, computers, automobiles, telecommunications and military applications. The lead frame electrically interconnects the integrated circuit chip to circuitry external to the device. The lead frame is typically formed from a highly electrically and thermally conductive material, such as copper or copper alloys. The lead frame material is stamped or etched into a plurality of leads, and a central area, called a die attach pad, onto which the integrated circuit chip is attached. The chip is electrically connected to the leads, usually by wire bonding, and the device is encapsulated to provide mechanical and environmental protection.\nConventional lead frames for integrated circuit packages are typically made using metal etching processes to create die attach pads and leads from a conductive layer of material on a bulk substrate or other material. Such process involves etching the conductive layer of material to form discrete conductor structures (e.g., the die attach pads and the leads). The etching process typically limits the spacing between the die attach pads and/or between adjacent leads on these conventional lead frames to at least one thickness of the conductive layer of material due to physical limitations of etching technology. In addition, the portions of the conductive layer removed from the substrate during the etching process generate waste in the manufacturing process. Manufacturer lead times for producing a new lead frame configuration via etching or other processes also can be substantial."}
{"text": "The invention relates to a mist separator apparatus that removes mist and liquid from structures containing the same.\nMany types of equipment that cut or otherwise form workpieces fabricated from hard material such as metal, ceramics or plastics sometimes require a cooling fluid to cool both the cutting tool and the workpiece as it is being cut or formed. One type of equipment is a wafer slicing machine 2 as schematically illustrated in FIG. 1 for purposes of explanation only. Generally, the wafer slicing machine 2 includes a housing 4 and a cutting assembly 6. A saw blade 8 connected to a shaft 10 is enclosed within the housing 4 and a motor 12 disposed exteriorly of the housing 4 drives the saw blade 8 to cut workpiece 14. Because a significant amount of heat is generated when cutting the workpiece 14 with the rotating saw blade 8, a liquid coolant 18, such as water or oil, is sprayed onto the rotating saw blade 8 and the workpiece 14 through a nozzle 20.\nMuch of the liquid coolant 18 which now contains minute particles of the workpiece 14 accumulates at the bottom of the housing 4. A mist M is also generated as a result of spraying the liquid coolant 18 onto the rotating saw blade 8 and the workpiece 14. For the wafer slicing machine 2 to operate efficiently, the liquid coolant 18 that now contains minute particles of the workpiece 14 must be drained from the housing 4 and the mist, some of which also contains minute particles of the workpiece 14, must be exhausted from the housing 4.\nA conventional drain/exhaust system 22 is illustrated in FIG. 1. A header 24 is attached exteriorly near the bottom of the housing 4 so that accumulated liquid coolant containing minute particles of the workpiece 14 can drain from the housing 4 through a drain pipe 26. Also, the mist M migrates from the housing 4 through the header 24 and into a leading exhaust pipe 28 into a canister 30. While in the canister 30, the heavier mist droplets tend to settle toward the bottom of the canister 30 while the lighter mist droplets are exhausted through a trailing exhaust pipe 32. As commonly known in the art, a slight vacuum is applied to the trailing exhaust pipe 32 to draw the lighter mist droplets therethrough.\nAs the heavier mist droplets accumulate at the bottom of the canister 30, a sludge-like material forms. If sufficiently wet, the sludge flows through a canister drain pipe 34 for disposal.\nUnfortunately, if excessive mist droplets are being exhausted through the trailing exhaust pipe 32, the sludge-like material tends to thicken and, as a result, does not effectively drain from the canister 30. Thus, a door 36 on the canister 30 must be opened to gain access into the canister 30 so it can be purged of the stagnate sludge-type material. Cleaning the canister 30 is messy and labor-intensive. Ironically, excessive mist droplets being exhausted through the trailing exhaust pipe 32 is known to be detrimental to the exhaust portion of the drain/exhaust system 22.\nThere is a need in the industry to provide a mist separator that removes a mist and a liquid material more efficiently and effectively from cutting or forming equipment. It would advantageous if the mist separator apparatus eliminates the use of the canister in order to eliminate the need to periodically clean the sludge-like material therefrom. It would also be advantageous if the mist separator apparatus can simultaneously reduce the amount of mist being exhausted through the vacuum-driven exhaust system while simultaneously increasing the amount of mist precipitating as liquids so the liquids can drain through the drainage system. There is also a need to provide a mist separator apparatus that inhibits the mist from directly entering into the drain/exhaust system. It would be advantageous to provide a mist separator apparatus that allows the liquid to initially drain into a drainage system and allow the mist to precipitate in a mist collecting container system which is initially separate and distinct from the drainage system. It would also be advantageous to provide a mist separator apparatus that directs the precipitating mist in the mist collecting container to merge downstream with the draining liquids. The present invention satisfies these needs and provides these advantages."}
{"text": "1. Field of the Invention\nThis invention relates to a pressure sensor and particularly to an inexpensive and simple pressure sensor having good static pressure characteristics.\n2. Description of the Related Art\nA conventional pressure sensor has a silicon base 9 and a boss 10 (see, for example, JP-UM-A-5-50335). The conventional pressure sensor will now be described in detail with reference to FIG. 1. FIG. 1 is a sectional view of the conventional pressure sensor.\nIn FIG. 1, a metal 5, which is a metal base, has a hole 6 to which a pressure is applied. Also a glass 4, which is a base, has the hole 6 to which a pressure is applied. The glass 4 and the metal 5 are mounted to each other by Au eutectic bond or the like. The glass 4 performs electrical and mechanical insulation.\nAlso, the silicon base 9 has the hole 6 to which a pressure is applied. The silicon base 9 and the glass 4 are mounted to each other. The boss 10 is provided between the silicon base 9 and the glass 4.\nA silicon sensor 1, which is a sensor, has a diaphragm 2 connected to the hole 6. The silicon sensor 1 has a strain gauge 3 for converting a strain (displacement) occurring in the diaphragm 2 to an electric signal. One side of the silicon sensor 1 is mounted to the silicon base 9.\nThe other side of the silicon sensor 1 contacts a room B. The hole 6 and the diaphragm 2 form a room A.\nA static pressure is applied to the silicon sensor 1, the silicon base 9, the glass 4 and the other parts in the conventional example of FIG. 1.\nIn the conventional example of FIG. 1 constructed as described above, the pressure applied to the hole 6 is converted to an electric signal. The strain gauge 3 in the conventional example of FIG. 1 generates an electric signal based on the differential pressure between the room A and the room B and the static pressure.\nWhen the static pressure is applied, the silicon sensor 1, the silicon base 9 and the glass 4 are deformed, respectively. Since the silicon sensor 1 and the silicon base 9 have a large Young's modulus, these are deformed slightly. Since the glass 4 has a small Young's modulus, it is deformed largely.\nThe silicon base 9 and the boss 10 in the conventional example of FIG. 1 restrain transfer of the influence of deformation of the glass 4 to the strain gauge 3 of the silicon sensor 1. The boss 10 reduces the bond area between the silicon base 9 and the glass 4.\nMoreover, some conventional pressure sensors (semiconductor pressure converting apparatus) separately have a structure for detecting a differential pressure and a structure for detecting a static pressure, and also have a structure for reducing interference with a differential pressure signal while increasing output of a static pressure signal (see, for example, Japanese Patent No. 2,656,566).\nHowever, the conventional example of FIG. 1 has a problem that an error occurs in the application of the static pressure (static pressure characteristics).\nSpecifically, since the Young's modulus of silicon is different from the Young's modulus of glass, when the static pressure is applied, the deformation of the silicon sensor 1 and the silicon base 9 differs from the deformation of the glass 4, and a strain based on the deformation of the glass is generated in the strain gauge 3.\nMoreover, since the glass 4 has characteristics such as delayed elasticity and viscoelasticity, it causes a strain in the diaphragm 2 and thus causes a strain in the strain gauge 3. This causes an error in the conventional example of FIG. 1.\nAlso, the silicon base 9 and the boss 10 in the conventional example of FIG. 1 have problems of increase in the number of components, increase in the number of processing steps, and high cost.\nMoreover, the formation of the boss 10 has a problem of deteriorating the yield of bond. The formation of the boss 10 also has a problem of lowering broken pressure reduce the bond area.\nMeanwhile, the conventional example of Japanese Patent No. 2,656,566 has a problem that it does not restrain occurrence of an error in the static pressure characteristics and cannot acquire good static pressure characteristics."}
{"text": "The present invention generally pertains to communications systems and is particularly directed to an improved base station for use in a subscriber communication network, such as a wireless digital telephone system.\nA base station for a wireless telephone system having a plurality of subscriber stations is described in U.S. Pat. No. 4,675,863 issued June 23, 1987, for \"Subscriber RF Telephone System For Providing Multiple Speech And/Or Data Signals Simultaneously Over Either A Single Or A Plurality of RF Channels\".\nThe base station described in the above-referenced patent communicates signals between subscriber stations and an external communication network having a plurality of ports. Such base station includes a communication circuit for enabling simultaneous communications between a plurality of the ports and a plurality of subscriber stations over a given communication channel having multiple sequentially repetitive time slots, with predetermined time slots being assigned respectively to predetermined subscriber stations; a remote-connection processor (RPU) for directing communications between the time slot assigned to a given subscriber station and a given port; and an exchange for connecting the communication circuit to the ports. The exchange includes a switch, which responds to a control signal from the RPU by physically connecting a selected port to a selected communication channel time slot assigned to a given subscriber station. The communication circuit includes a plurality of channel control units (CCUs) for coupling the assigned communication channel time slot to the given subscriber station in response to a command signal communicated by the RPU to CCU's over a base-station control channel (BCC) in response to status messages provided over the BCC to the RPU to indicate the usage status of the communication channel time slots and the subscriber stations. The assigned communication channel time slot is coupled to the given subscriber station by an assigned time slot in an assigned radio frequency (RF) channel. The BCC is provided over lines seprately connected directly from the RPU to each of the CCUs. Control commands and status messages are communicated between the CCUs and the subscriber stations over a radio control channel (RCC) assigned to a predetermined time slot of a predetermined RF channel."}
{"text": "In situations where power lines exist between two devices which wish to communicate to each other, it has been recognised that the existing power lines may be used both for power transmission and communications, thereby avoiding the need to install separate communication cabling. In the past, this has been achieved by way of high speed power line modems, which are capable of connecting to the power line and inserting a carrier signal onto the power lines and additional to the power signal, and modulating the carrier signal with information to be transmitted. However, the circuitry of such modems must be able to withstand the large power levels present on the power lines, and the cost of such modems is usually significantly higher than standard modems which can be used where dedicated communications cabling is available, thereby offsetting some of the cost savings in using only a single line for both power and communications. Furthermore, devices connected to the power line are often quite simple in nature and consequently the required rate of communications will be extremely low, as only relatively small amounts of information need to be transmitted. In such situations, the high speed and high cost of a power line modem will go substantially unutilised.\nAny discussion of documents, acts, materials, devices, articles or the like which has been included in the present specification is solely for the purpose of providing a context for the present invention. It is not to be taken as an admission that any or all of these matters form part of the prior art base or were common general knowledge in the field relevant to the present invention as it existed before the priority date of each claim of this application.\nThroughout this specification the word “comprise”, or variations such as “comprises” or “comprising”, will be understood to imply the inclusion of a stated element, integer or step, or group of elements, integers or steps, but not the exclusion of any other element, integer or step, or group of elements, integers or steps."}
{"text": "Butyric acid (BA) is a natural product. It is supplied to mammals from two main sources: 1) the diet, mainly from dairy fat, and 2) from the bacterial fermentation of unabsorbed carbohydrates in the colon, where it reaches mM concentrations (Cummings, Gut 22:763-779, 1982; Leder et al., Cell 5:319-322, 1975).\nBA has been known for nearly the last three decades to be a potent differentiating and antiproliferative agent in a wide spectra of neoplastic cells in vitro (Prasad, Life Sci. 27:1351-1358, 1980). In cancer cells, BA has been reported to induce cellular and biochemical changes, e.g., in cell morphology, enzyme activity, receptor expression and cell-surface antigens (Nordenberg et al., Exp. Cell Res. 162:77-85, 1986; Nordenberg et al., Br. J. Cancer 56:493-497, 1987; and Fishman et al., J. Biol. Chem. 254:4342-4344, 1979).\nAlthough BA or its sodium salt (sodium butyrate, SB) has been the subject of numerous studies, its mode of action is unclear. The most specific effect of butyric acid is inhibition of nuclear deacetylase(s), resulting in hyperacetylation of histones H3 and H4 (Riggs, et al., Nature 263:462-464, 1977). Increased histone acetylation following treatment with BA has been correlated with changes in transcriptional activity and the differentiated state of cells (Thorne et al., Eur. J. Biochem. 193:701-713, 1990). BA also exerts other nuclear actions, including modifications in the extent of phosphorylation (Boffa et al., J. Biol. Chem. 256:9612-9621, 1981) and methylation (Haan et al., Cancer Res. 46:713-716, 1986). Other cellular organelles, e.g., cytoskeleton and membrane composition and function, have been shown to be affected by BA (Bourgeade et al., J. Interferon Res. 1:323-332, 1981). Modulations in the expression of oncogenes and suppressor genes by BA were demonstrated in several cell types. Toscani et al., reported alterations in c-myc, p53 thymidine kinase, c-fos and AP2 in 3T3 fibroblasts (Oncogene Res. 3:223-238, 1988). A decrease in the expression of c-myc and H-ras oncogenes in B16 melanoma and in c-myc in HL-60 promyelocytic leukemia was also reported (Prasad et al., Biochem. Cell Biol. 68:1250-1255, 1992; and Rabizadeh et al., FEBS Lett. 328:225-229, 1993).\nBA has been reported to induce apoptosis, i.e., programmed cell death. SB has been shown to produce apoptosis in vitro in human colon carcinoma, leukemia and retinoblastoma cell lines (Bhatia et al., Cell Growth Diff. 6:937-944, 1995; Conway et al., Oncol. Res. 7:289-297, 1993; Hague et al.; Int. J. Cancer 60:400-406, 1995). Apoptosis is the physiological mechanism for the elimination of cells in a controlled and timely manner. Organisms maintain a delicate balance between cell proliferation and cell death, which when disrupted can tip the balance between cancer, in the case of over accumulation of cells, and degenerative diseases, in the case of premature cell losses. Hence, inhibition of apoptosis can contribute to tumor growth and promote progression of neoplastic conditions.\nThe promising in vitro antitumor effects of BA and BA salts led to the initiation of clinical trials for the treatment of cancer patients with observed minimal or transient efficacy. [Novogrodsky et al., Cancer 51:9-14, 1983; Rephaeli et al., Intl. J. Oncol. 4:1387-1391, 1994; Miller et al., Eur. J. Cancer Clin. Oncol. 23:1283-1287, 1987].\nClinical trials have been conducted for the treatment of .beta.-globin disorders (e.g., .beta.-thalassemia and sickle-cell anemia) using BA salts. The BA salts elevated expression of fetal hemoglobin (HbF), normally repressed in adults, and favorably modified the disease symptoms in these patients (Stamatoyannopouos et al., Ann. Rev. Med. 43:497-521, 1992). In this regard, arginine butyrate (AB) has been used in clinical trials with moderate efficacy (Perrine et al., N. Eng. J. Med. 328:81-86, 1993; Sher et al., N. Eng. J. Med. 332:1606-1610, 1995). The reported side effects of AB included hypokalemia, headache, nausea and vomiting in .beta.-thalassemia and sickle-cell anemia patients.\nButyric acid derivatives with antitumor activity and immunomodulatory properties have been reported in U.S. Pat. No. 5,200,553 and by Nudelman et al., 1992, J. Med. Chem. 35:687-694. The most active butyric acid prodrug reported in these references was pivaloyloxymethyl butyrate (AN-9). None of the compounds disclosed in these references included carboxylic acid-containing oxyalkyl compounds of this invention.\nBA and/or its analogues have also been reported to increase the expression of transfected DNA (Carstea et al., 1993, Biophys. Biohem. Res. Comm. 192:649; Cheng et al., 1995, Am. J. Physical 268:L615-L624) and to induce tumor-restricted gene expression by adenovirus vectors (Tang et al., 1994, Cancer Gene Therapy 1:15-20). Tributyrin has been reported to enhance the expression of a reporter gene in primary and immortalized cell lines (Smith et al., 1995, Biotechniques 18:852-835).\nHowever, BA and its salts are normally metabolized rapidly and have very short half-lives in vivo, thus the achievement and maintenance of effective plasma concentrations are problems associated with BA and BA salts, particularly for in vivo uses. BA and BA salts have required large doses to achieve even minimal therapeutic effects. Because of the high dosage, fluid overload and mild alkalosis may occur. Patients receiving BA emanate an unpleasant odor that is socially unacceptable.\nWhile BA salts have been shown to increase HbF expression, and appear to hold therapeutic promise with low toxicity in cancer patients, they nevertheless have shown low potency in in vitro assays and clinical trials. There also remains a need to identify compounds as effective or more effective than BA or BA salts as differentiating or anti-proliferating agents for the treatment of cancers. Such compounds need to have higher potency than BA without the problems associated with BA (such as bad odor). Consequently, there remains a need for therapeutic compounds that either deliver BA to cells in a longer acting form or which have similar activity as BA but a longer duration of effectiveness in vivo.\nThe compounds of this invention address these needs and are more potent than BA or BA salts for treating cancers and other proliferative diseases, for treating gastrointestinal disorders, for wound healing and for treating blood disorders such as thalassemia, sickle cell anemia and other anemias, for modulating an immune response, for enhancing recombinant gene expression, for treating insulin-dependent patients, for treating cystic fibrosis patients, for inhibiting telomerase activity, for detecting cancerous or malignant cells, for treating virus-associated tumors, especially EBV-associated tumors, for augmenting expression of a tumor suppressor gene and for inducing tolerance to an antigen. For example, one of the advantages of the compounds of the invention is increased water solubility of the free carboxylic acids compounds of the invention and their salts, and easier administration, especially for intravenous administration."}
{"text": "Many computer systems are used in fault tolerant environments such as the banking industry, airline industry, mercantile markets, and so forth. In such environments any downtime to a computer system caused by a computer system failure can be catastrophic to businesses and customers that rely on such systems. To safeguard against such failures, many computer systems rely on maintenance systems to monitor the health of the systems and take corrective action in the event problems are detected.\nUnfortunately, many of these maintenance systems are unable to precisely diagnose the root cause of a problem or take corrective action fast enough to prevent the computer system from failing. To compound problems further, many maintenance systems monitor the health of computer systems by testing them in a “service mode.” Typically, when a computer is placed into a service mode, the normal operational mode of the computer system is interrupted or degraded. Additionally, the service mode creates an artificial testing atmosphere in which problems that may occur in the normal mode of operation never present themselves in the service mode, and hence, may go undetected, possibly jeopardizing the operating stability of the computer system."}
{"text": "Traditionally firearm barrels have been manufactured out of steel alone, but users have long desired lighter weight firearms that remain durable and reliably accurate. In order to address this desire, it has previously been known to construct strong, lightweight barrels for firearms using composite materials. For example, some previously known composite gun barrels include an inner tubular structure made from a hard material such as steel alloy and an outer jacket made from a composite material such as a continuous carbon fiber-reinforced polymer matrix composite. This combination lightens the gun while retaining barrel strength and stiffness.\nThe manner in which these previously known gun barrels have been made includes applying carbon fiber in a wet filament winding operation, wherein dry carbon fiber strands are combined with a resin (such as a crosslinkable epoxy, a polyimide, cyanate ester, inorganic polymer or thermoplastic polymer) in a wet dip pan process, then wound around an inner liner and processed. The composite barrel may then be cured such that the resin bonds the outer shell to the inner liner, followed by attaching the barrel to a receiver and stock.\nSuch carbon fiber-reinforced composites may provide a suitable balance of thermal properties, mechanical properties, and processing characteristics for many common firearms applications. However, such composite gun barrels can also pose problems not encountered with traditional steel barrels. For example, difficulty may arise in constructing the composite barrel in a manner and quantity around and along the liner that ensures that the barrel does not burst upon firing, achieving satisfactory strength and stiffness in the principal directions (for example, axially and torsionally), providing adequate environmental durability, and dampening the shock wave that propagates when a projectile is fired.\nSome of the above issues can be addressed by using additional windings (i.e., more circumferential “hoop wraps”) to improve burst strength and more axially oriented helical windings to improve axial tensile and flexural strength and stiffness. However, adding more layers of windings can lead to manufacturing and curing complications, higher material expense, more weight, and/or a bulkier barrel profile than desired.\nThermal management is also a significant problem when using continuous fiber composite (CFC) material for the outer shells because the resin content within the CFC is a relatively poor conductor (i.e., insulator of the heat generated by hot gasses within the liner). Additional layers of CFC windings exacerbate the heat removal problem. During operation, the barrel will heat up. In the case where the matrix phase is an organic polymer, if the cured resin within the CFC reaches its glass transition temperature, Tg, the CFC softens significantly and the mechanical integrity of the composite barrel is compromised. As the barrel is heated to even higher temperatures, irreversible thermal decomposition of the cured matrix occurs and barrel structural integrity is further compromised.\nA further problem relates to stresses within the barrel arising from thermal expansion differences between the composite outer liner and the inner liner of the composite barrel. As the inner steel liner heats during operation, it expands both radially and longitudinally. Composite structures in the prior art have a substantially lower average effective coefficient of thermal expansion (CTE) in the longitudinal direction than steel and so when heated, the CFC outer shell expands substantially less than the steel liner. This may increase or decrease thermal stresses in the barrel depending on the state of thermal residual stress from processing. For example, when a stainless steel liner and a typical CFC are subjected to heating during operation, uneven expansion can produce thermal stresses on the liner-CFC interface, possibly even causing separation of the CFC from portions of the liner or fractures within the CFC shell. Even if no separation occurs, minor variations in the CFC or metal liner properties, or geometric variation, may promote uneven thermal stresses at the interface between the barrel and CFC that may result in nonlinear deformation or displacement of the barrel from its original axis. Even a very slight displacement can significantly degrade accuracy. Furthermore, even if the barrel and liner remain perfectly true, the various layers of windings within the CFC can have different CTEs, especially longitudinally. When subjected to elevated operating temperatures, differences in the thermal expansion of adjacent winding layers within the CFC can result in high level of interlaminar shear stress and even delamination.\nFrom the above, it can be seen that there is a need to produce an improved composite barrel that uses materials that provide superior axial and torsional strength and stiffness while minimizing weight, radial bulk and interlaminar stress, resists deformation when heated and reduces high amplitude vibrations at the muzzle end of the barrel."}
{"text": "1. Field\nOne or more embodiments disclosed herein relate to a scan driving circuit and a driving method for the scan driving circuit.\n2. Description of the Related Art\nPersonal computers, portable phones, portable information terminals, and monitors of various information devices are all equipped with displays. The display may be, for example, a liquid crystal display, an organic light emitting diode display, or a plasma display panel. An organic light emitting diode display in particular has excellent emission efficiency, luminance, viewing angle, and response time.\nAn organic light emitting diode display generates images using pixels that include organic light emitting diodes. The pixels are arranged in a matrix at cross points of data lines, scan lines, and power supply lines. In addition to an organic light emitting diode, each pixel includes a driving transistor and one or more capacitors. Light is emitted from the diode based on a recombination of electrons and holes in a light emitting active layer.\nOver time, the amount of current flowing in the organic light emitting diode of each pixel may change based on a deviation in the threshold voltage of the driving transistor. As a result, non-uniformity may occur in the display. Additionally, the characteristics of the driving transistor may change based on manufacturing process parameters. Because it is difficult for the transistors in an organic light emitting diode display to have the same characteristics, the threshold voltage deviation of the driving transistors in the pixels may be different.\nIn attempt to overcome this deviation, a compensation circuit may be used in each pixel. The compensation circuit may charge a voltage corresponding to the threshold voltage of the driving transistor for 1 horizontal period. However, in some displays and especially in a high-resolution organic light emitting diode display, coupling may occur in a parasitic capacitor between lines and horizontal striped patterns. This may adversely affect display performance."}
{"text": "1. Field of the Invention\nThis invention relates to austenitic stainless steels and nickel-base alloys, particularly such alloys, and methods of making the same, wherein the alloys are strengthened by nanometer-size hollow oxides which serve as nucleation sites for chromium-rich carbide precipitates within the alloy grains.\n2. Prior Art\nStrengthening of metallic alloys primarily is achieved through alloy grain size control, solute additions to a base metal to produce solid solution strengthening, and/or dispersion (precipitation or second phase) strengthening effects. These methods have been applied to a variety of metallic alloy systems and are the basis for strengthening of many of the high-value alloys available in the metal market today. For example, according to U.S. Pat. No. 4,758,405, high strength Al alloys have been produced by gas atomization of an Al alloy melt with an inert gas such as argon, helium or nitrogen containing 0.5-2% by volume of oxygen. U.S. Pat. No. 4,999,052 discloses austenitic stainless steels strengthened with nitrogen in solid solution and containing a dispersant such as a nitride, for example, titanium nitride, and/or an oxide such as yttria. The role of nitrogen in iron-base alloys, particularly austenitic stainless steels, has received considerable attention during the past 80 years. Two fairly recent symposia on this subject have provided state-of-the-art reviews. Proceedings of the International Conference on High-Nitrogen Steels-88, Editors J. Foct and A Hendry, Publ. Institute of Metals, London GB (1988); Proceedings of the International Conference on High-Nitrogen Steels-90, Editors G. Stein and H. Witulski, Publ. Verlag Stalil Eisen, MbH, Dusseldorf (1990).\nAlloy 654SMO is a relatively new austenitic stainless steel of high strength and good corrosion resistance. B. Wallen, M. Liljas and P. Stenvall, Avesta 654 SMO--a New High Molybdenum, High Nitrogen Stainless Steel, Avesta Corrosion Management, Avesta AB, S-774 80 Avesta, Sweden.\nAn overview of mechanisms for strengthening austenitic stainless steels is provided by K. J. Irvine et al., \"High-Strength Austenitic Stainless Steels,\" Journal of The Iron and Steel Institute, October 1961.\nAlloys 625 and 718 are representative of high strength nickel-base alloys. H. L. Eiselstein et al. \"The Invention and Definition of Alloy 625,\" Inco Alloys International, Inc, P.O. Box 1958, Huntington, W. Va., Superalloys 718, 625 and Various Derivatives, E. A Loria, Ed., The Minerals, Metals & Materials Society, 1991. Such alloys have been produced by the powder metallurgy process. F. J. Rizzo et al. \"Microstructural Characterization of PM 625-Type Materials,\" Crucible Compaction Metals, McKee and Robb Hill Roads, Oakdale, Pa. 15071 and Purdue University, West Lafayette, Ind. 47906, included in Superalloys 718, 625 and Various Derivatives, E. A Loria, Ed., The Minerals, Metals & Materials Society, 1991. See also R. B. Frank \"Custom Age 625 Plus Alloy--A Higher Strength Alternative to Alloy 625, Carpenter Technology Corporation, P.O. Box 14662, Reading, Pa. 19612, also included in Superalloys 718, 625 and Various Derivatives, E. A Loria, Ed., The Minerals, Metals & Materials Society, 1991.\nHowever, the options for improving the properties and performance of metallic alloys are becoming limited in terms of new developments, and new, innovative methods are needed in order to provide a new generation of advanced alloys that can stand up to increasingly severe future demands."}
{"text": "Functional and parametric degradation of microcircuits due to total ionizing dose (TID) creates significant obstacles to the deployment of critical state of the-art technologies in space missions. Moreover, since device dielectrics in which such degradation occurs vary from one fabrication lot to the next, these effects must be re-evaluated on a lot-by-lot basis. Space mission designers are moving toward smaller, cheaper designs that reduce system-level shielding and consequently, limit the lifetimes of the space missions. Currently, mission radiation reliability requirements are achieved by shielding the electronics box and/or by the shielding on the satellite or spacecraft in conjunction with the use of radiation-hardened electronics. However, space missions are targeting more radiation-harsh environments thereby demanding improved shielding techniques. For example, vulnerable components typically require shielding over 4 pi-steradians. Often, the most effective mitigation against total ionizing dose (TID) degradation is the addition of radiation shielding to the electronics box. Unfortunately, shielding materials can add significant mass to a system. One conventional method for reducing additional mass to the system is to apply spot shielding in the form of cover plates only on the critical components that require shielding. This conventional method provides limited shielding because of the omnidirectional natural radiation environment in space. However, reducing radiation shielding for the electronics box and/or reducing spacecraft-level shielding will necessitate more complex spot shielding to protect the components from the omnidirectional radiation environment. Such complex spot shielding entails complex patterned designs that require labor-intensive machining of complex metal designs and time-consuming installation. Furthermore, the shielding effectiveness of such complex patterned designs must be verified with three-dimensional ray-trace analysis. Small package sizes make component-level shielding fabrication even more challenging.\nWhat is needed is a method for forming a component-level radiation shield that eliminates the problems and disadvantages associated with the aforementioned conventional practices."}
{"text": "Photodynamic therapy (PDT), which is also called as photochemotherapy, is a technique for treating incurable disease such as cancer without operation or treating skin disease such as acne by the use of a photosensitizer. There has been active study about PDT since the early twentieth century. Currently, PDT is used to increase immunocompetence in diagnosis and treatment of cancer, autologous bone marrow transplantation, antibiotics, AIDS treatment, skin graft or treatment of arthritis and the like, and so its applications have been gradually expanded. Specifically, PDT used in the treatment of cancer is a therapeutic method utilizing a principle in which a photosensitizer—which is a chemical compound showing sensitivity to light—is administered to the body; when it is exposed to external light, singlet oxygen or free radicals are generated via the chemical reaction by rich oxygen in the body and external light; and then such singlet oxygen or free radicals destroy various lesion sites or cancer cells by inducing cell death.\nAs photosensitizers used in PDT, porphyrin derivatives, chlorin, bacteriochlorin, phthalocyanine, 5-aminolevulinic acid derivatives and the like are known. Cyclic tetrapyrrole derivatives as photosensitizers also have properties for being utilized as agents for early diagnosis of tumor since they show fluorescence and phosphorescence according to their chemical properties as well as being selectively accumulated to cancer cells. In addition, metalloporphyrin, in which metal is bound to the interior of cyclic tetrapyrrole, shows various properties according to the kind of bound metal, thereby being applied to early diagnosis of tumor cells such as cancer cells by use of a diagnostic technique wherein metalloporphyrin is used as an MRI (magnetic resonance imaging) contrasting agent. 5-Aminolevulinic acid derivatives, which are the most widely known photosensitizers, have advantages wherein the method of use is simple, they can relatively easily penetrate into the skin due to small molecular weight, and they are safe with few side effects. Furthermore, one of the photosensitizers receiving attention is a chlorin-type photosensitizer, chlorin e6. As an example of the composition using chlorin e6, Korean Patent Application No. 10-2006-0136610 discloses antitumoric compositions for oral administration comprising chlorin e6 or a pharmaceutically acceptable salt thereof as an active ingredient.\nAs explained above, usage in the treatment of skin diseases such as acne was suggested as one of the applications of PDT, but no report has been made about its successful application."}
{"text": "Conventionally, as disclosed in Patent Literature 1, an on-board electronic device having CPUs and substrates, respectively corresponding to types of functions, such as a navigation board on which a CPU to perform navigation processing is mounted and an audio board on which a CPU to perform audio processing is mounted, is known. Note that this type of on-board electronic device is often integrated with a display and a touch panel.\nFurther, in recent years, in accordance with multimedization in the vehicle, a configuration, where a user interface unit (UI unit) having integrated display and touch panel is connected to a plurality of electronic devices mounted on the vehicle such as a navigational device and an audio instrument, via a display control unit (hereinbelow, DCU), is increasingly adopted. The DCU plays a role in control of the respective operations of the plurality of electronic devices connected to the own device, and generation (or acquisition) of image data to be displayed on the display based on requests from the respective electronic devices, and display of the image data on the display. Hereinbelow, a unit where the UI unit and the DCU are integrated will be referred to as a display unit for vehicle.\nA conventional DCU has two substrates, i.e., a main board mounted with a main CPU and a sub board mounted with a sub CPU and a power supply circuit. The main CPU mainly controls the operations of the electronic devices connected to the own device, and draws an image to be outputted to the display. The sub CPU performs management (acquisition and storage) of vehicle information inputted from an in-vehicle network, and controls electric power supply to the respective elements of the display unit for vehicle in cooperation with the power supply circuit.\nThe conventional DCU has two substrates, i.e. the main board and the sub board, mainly for reducing the area of each substrate to an area smaller than that of the display.\nMore particularly, to realize a function, performed with a main board and a sub board, using one substrate, the area of the substrate is larger than the area of the display. When the area of the substrate is larger than the area of the display, its housing is also larger than the display, to impair mountability to the vehicle.\nAccordingly, in the conventional DCU, a function to be provided in the DCU is shared with the two substrates, the main board and the sub board, so that the size of each substrate is suppressed to be equal to or smaller than the size of the display.\nHowever, in the conventional display unit for vehicle (more particularly, DCU) having two substrates, the manufacturing cost such as costs of parts and machining cost is increased. The machining cost here includes costs necessary for assembling work of the respective substrates in the housing and wiring connection work between the substrates.\nFurther, it is desired to increase the CPU clock for improvement in DCU performance. However, in the conventional DCU, the space in the housing is partitioned with the two substrates. When the clock is increased, radiation with a cooling fan becomes insufficient. Accordingly, from the viewpoints of manufacturing cost and heat dissipation, it is preferable that the number of the substrates is one.\nOn the other hand, as described above, when a function realized with two substrates is simply integrated in one substrate, the area of the substrate becomes larger than the area of the display, which might impair the mountability. Especially, when the difference between the area of the display and the area of the substrate is large, i.e., when the area of the substrate is larger, the mountability is impaired."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for recording and/or reproducing information on an optical recording medium, for example a magneto-optical disk drive, a direct read after write drive (DRAW), a phase changing disk drive, CD-ROM or the like.\n2. Related Art Statement\nThe disks for use in such an apparatus are provided with a protective layer (hereinafter, referred to as a cover glass) such as glass, resin (PC, PMMA or the like) formed on its side surface on which the reading light and the writing light are incident, in order to protect or guard the recording layer of the recording medium. This cover glass has different thicknesses in accordance with the types of the disk, even in the same DRAW disk. For example, the thicknesses of the disks are 1.2 mm for the magneto-optical disk, 1.2 mm and 0.6 mm for the phase changing disk, and 1.2 mm and 0.09 mm for the standard DRAW disk of 356 mm.\nAn objective lens for use in an optical pick-up of the information recording and/or reproducing apparatus has generally a number aperture NA of 0.45-0.6, and is designed by taking aberration caused by the above-discussed cover glass into consideration. On considering allowable level at the reading/writing property of the signal, in the case of the cover glass having a 1.2 mm thickness, the limit of the allowable level is in the order of xc2x10.05. If the allowable level exceeds this limit value, the reading/writing property of the signal becomes considerably deteriorated. If the reading/writing operation of the signal is performed by exchanging the disks having a different thickness of the cover glass, for example, the disk having a 1.2 mm thickness of the cover glass and the disk having a 0.6 mm thickness of the cover glass, therefore, it is difficult to perform this reading/writing operation of the signal with one information recording and/or reproducing apparatus. That is, two information recording and/or reproducing apparatus must be required.\nIn order to remove the above defect, for example, Japanese Patent Application Opened No. 241,095/93 provides the feature of inserting parallel plates between a light source and a collimator to compensate a spherical aberration caused by the difference in thickness of the cover glass of the disk.\nHowever, the optical pick-up disclosed in the Japanese Patent Application Opened No. 241,095/93 has a problem as follows:\nFor example, a spherical aberration component xcfx8940 of a coefficient of a wave front aberration and a value of the wave front aberration, which are caused, for example, in case of passing the light through the cover glass having thickness t can be obtained from xe2x80x9cOptics 14xe2x80x9d (1985), pp. 219-221 and expressed as follows:\nxcfx8940=(t/8)xe2x88x92{(n2xe2x88x921)/n3}xe2x88x92NA4xe2x80x83xe2x80x83(1)\n\"ugr\"=0.0745|xcfx8940|xe2x80x83xe2x80x83(2)\nwherein the numeral aperture of the objective lens is NA, and the refractive index is n. Therefore, in the pick-up having NA=0.55 and n=1.57 and using the light of wavelength xcex=780 nm, if the disk having the cover glass of t=1.2 mm is exchanged to the disk having the cover glass of t=0.6 mm, the spherical aberration of xcfx8940=0.00260, \"ugr\"=0.248 xcexrms is generated.\nWhile in case of compensating the above-discussed spherical aberration by arranging the parallel plates for compensation between the light source and the collimator lens, if the numerical aperture at the light source side is NAxe2x80x2=0.25 and the refractive index is n=1.57, the thickness txe2x80x2 of the parallel plates is expressed as follows:\ntxe2x80x2={n3/(n2xe2x88x921)}xe2x88x92(8/NAxe2x80x24)xe2x88x92xcfx8940=14.07 mm\nMoreover, in order to insert such thick parallel plates between the light source and the collimator lens, the light source must be moved by a following distance 1:\nl=txe2x80x2(1xe2x88x921/n)=5.11 mm.\nIn such a way, the optical pick-up disclosed in the Japanese Patent Application Opened No. 241,095/93 requires a moving mechanism for inserting the compensating parallel plates and a mechanism for moving the light source in accordance with the insertion of the parallel plates, so that the construction of the apparatus becomes complicated and the optical system becomes large. In the reading/writing operation, moreover, the objective lens is moved with the following up of the eccentricity of the disk by a tracking servo, thereby causing non-coincidence of the light axis of the objective lens to the light axis of the collimator lens, resulting in a generation of comma aberration or stigmatism.\nIt is an object of the present invention to eliminate the above described disadvantages of the conventional information recording and/or reproducing apparatus.\nIt is another object of the present invention to provide an information recording and/or reproducing apparatus capable of performing the reading/writing operation for plural kinds of optical recording media having cover glasses of a different thicknesses by means of a unity apparatus.\nAccording to a first aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source; at least two objective lenses corresponding to at least two kinds of optical recording media having different thicknesses of a protective layer, respectively; a holder for supporting these objective lenses; a selection means for selecting these objective lenses in accordance with the optical recording medium; the light of the light source being illuminated on the optical recording medium through the objective lens selected by the selection means, thereby performing a recording and/or reproducing of the information.\nAccording to a second aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source; an objective lens corresponding to an optical recording medium having a predetermined thickness of a protective layer; a holder for supporting the objective lens; an optical element for converting a wave front for an optical recording medium having a protective layer of thickness different from the above optical recording medium; a selection means for inserting the optical recording medium in an optical path between the objective lens and the optical recording medium, selectively and detachably, in accordance with the thickness of the protective layer of the optical recording medium to be recorded and/or reproduced for the information\nAccording to a third aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source an objective lens for generating at least two light spots corresponding to at least two kinds of optical recording media having different thicknesses of a protective layer; a holder for supporting these objective lenses; the light from the light source being illuminated on the optical recording medium through the objective lens, thereby performing recording and/or reproducing of the information.\nAccording to a fourth aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source an objective lens corresponding to an optical recording medium having a predetermined thickness of a protective layer; a holder for supporting the objective lenses; objective lens drive means for driving the objective lens relatively for the optical recording medium, a correction optical element capable of displacing the lens unitedly, for compensating a spherical aberration for an optical recording medium having the thickness different from the above optical recording medium.\nAccording to a fifth aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source; two optical systems corresponding to two kinds of optical recording media having different thicknesses of a protective layer; a selection means for selecting these optical systems in accordance with the optical recording medium; characterized by comprising: the selection means for detecting the information provided on a cartridge having the optical recording medium accommodated therein, and for selecting either one of said two optical systems.\nAccording to a sixth aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source; two optical systems corresponding to two kinds of optical recording media having different thicknesses of a protective layer; a selection means for selecting these optical systems in accordance with the optical recording medium; characterized by comprising: the selection means for detecting the position on the information recording surface of the optical recording medium and the thickness of the protective layer, and for selecting either one of said two optical systems.\nAccording to a seventh aspect of the present invention, there is provided an information recording and/or reproducing apparatus comprising: a light source; at least two optical systems corresponding to at least two kinds of optical recording media having different thicknesses of a protective layer; a selection means for selecting these optical systems in accordance with the optical recording medium; the light from the light source being illuminated on the optical recording medium, thereby performing a recording and/or reproducing of the information; and a detection means for detecting which objective lens of said at least two objective lenses.\nIn an embodiment of the information recording and/or reproducing apparatus according to the present invention, the selection means includes at least a part of a lens drive means for controlling a relative position of the objective lens for the optical recording medium by displaying the objective lens in predetermined direction during the recording and/or reproducing of the information.\nIn this case, the construction becomes simple and the apparatus can be downsized and inexpensive.\nIn a preferable embodiment of the information recording and/or reproducing apparatus according to the present invention, the correction optical element is so constructed that it is intervened selectively in the optical path between the optical recording medium and the objective lens with the use of at least a part of a lens drive means for controlling a relative position of the objective lens for the optical recording medium by displaying the objective lens in predetermined direction during the recording and/or reproducing of the information.\nTherefore, the construction becomes simple and the apparatus can be downsized and inexpensive.\nAccording to the first aspect of the present invention, the reading/writing operation can be performed for the plural kinds of the optical recording media in the single apparatus without causing the deterioration of the optical property due to the difference in thickness of the protective layer of the optical recording medium, by selecting the suitable objective lens with the selection means in accordance with kinds of the optical recording medium due to the thickness of the protective layer. The plural objective lenses are exchanged and selected in accordance with the kinds of the optical recording medium and the objective lens can be optionally selected as compared with the case of inserting the correction optical element between the objective lens and the optical recording medium, the objective lens having a small operating distance may be used, and the overall construction of the apparatus can be made thin.\nAccording to the second aspect of the present invention, the correction optical element capable of being displaced with the objective lens integrally is intervened selectively in the optical path between the optical recording medium and the objective lens in accordance with kinds of the optical recording medium due to the thickness of the protective layer, so that the reading/writing operation can be performed for the plural kinds of the optical recording media in the single apparatus without causing the deterioration of the optical property due to the difference in thickness of the protective layer of the optical recording medium. Moreover, the correction optical element is displaced with the objective lens integrally, so that they do not collide and thus the correction optical element can be positioned near the objective lens. Therefore, the objective lens having a small operating distance may be used, as compared with the case of inserting only the correction optical element always between the objective lens and the optical recording medium, and thus the overall construction of the apparatus can be made thin."}
{"text": "1. Field of the Invention\nThis invention relates to a fence support that is inexpensive and may be easily and quickly installed.\n2. Background Discussion\nTypical rural fences consist of a fencing material such as wired type mesh or barbed wire stretched over posts made of wood or metal. These posts are either driven into the ground with sledgehammers or powered or manual drivers, or a hole is dug into the ground, the post is inserted into the hole, and then concrete is poured into the hole and allowed to harden. Installing such fence posts is both laborious and time consuming For example, it normally takes about an hour to install a post by digging a hole, inserting the post, and pouring in the concrete.\nIt is sometimes desirable to brace the post using, for example, angle iron, a tee post, or tubing. One end of the brace is bolted to the top of the post and the other end is inserted into a hole in the ground and filled with concrete. A right triangle is formed, with the post and surface of the ground being the legs of triangle and the brace being the hypothenuse. Sometimes the braces pounded into the ground and then bolted or welded to the top of the post. This requires some skill and luck to make the brace intersect the top of the post at the desired position and angle. Frequently, the brace does not end up were desired.\nIf the post is not in the ground deep enough, or if the ground is soft, for example sandy or wet, the post is pulled from the ground when the fencing material is stretched from the post. The braces remain in place, but the post rises from the ground in much the same way that a pole vaulter is vaulted from the ground when he makes his jump. When this happens there is not a good solution. Typically, sand bags or rocks are placed around the post and positioned to exert a force in a direction opposite to the force being exerted by the stretched fencing material."}
{"text": "The ability to incorporate non-genetically encoded amino acids (i.e., “non-natural amino acids”) into proteins permits the introduction of chemical functional groups that could provide valuable alternatives to the naturally-occurring functional groups, such as the epsilon —NH2 of lysine, the sulfhydryl —SH of cysteine, the imino group of histidine, etc. Certain chemical functional groups are known to be inert to the functional groups found in the 20 common, genetically-encoded amino acids but react cleanly and efficiently to form stable linkages with functional groups that can be incorporated onto non-natural amino acids.\nMethods are now available to selectively introduce chemical functional groups that are not found in proteins, that are chemically inert to all of the functional groups found in the 20 common, genetically-encoded amino acids and that may be used to react efficiently and selectively with reagents comprising certain functional groups to form stable covalent linkages."}
{"text": "This invention generally relates to a control switch and device that includes tailored feedback to assist a person in adjusting various functions for a variety of systems.\nModern mass production vehicles tend to include more and more special features or subsystems. The addition of more features to a vehicle provides the potential for a more enjoyable driving experience; however, introducing more and more features or subsystems into a vehicle is not without problems.\nAn important need is to provide a driver the ability to control or adjust the various available subsystems without unduly distracting the driver from the task of driving the vehicle. It is important, therefore, to provide a driver with input control devices that allow a driver to realize that a desired adjustment or function selection has been accomplished without diverting the driver's attention from the road.\nTherefore, there is a need for a device that is readily usable by a driver to control a variety of the subsystems available within a vehicle. Prior to this invention, no satisfactory solution has been found that addresses the various needs discussed above. This invention is a device that provides customized feedback to a driver, which is indicative of the adjustment or selection that has been made. Moreover, this invention provides a device that is readily usable for controlling a variety of functions for a variety of vehicle subsystems."}
{"text": "A memory device having a three-dimensional structure has been proposed in which memory holes are made in a stacked body in which electrode layers that function as control gates of memory cells are multiply stacked with insulating layers interposed between the electrode layers, and silicon bodies used to form channels are provided on the side walls of the memory holes with a charge storage film interposed between the silicon bodies and the side walls.\nIn such a three-dimensional device, the memory holes and slits are made in the stacked body that includes the multiple layers of electrode layers and the multiple layers of insulating layers by, for example, RIE (Reactive Ion Etching). High shape controllability and dimensional controllability are desirable for the etching in such a case."}
{"text": "The present invention relates to firearms in general and in particular to a novel safety mechanism, a magazine-latching mechanism and a firing trigger mechanism in a 45 caliber recoil-operated blowback type semi-automatic pistol and the like.\nIn conventional automatic and semi-automatic firearms there is generally provided a magazine for holding a plurality of rounds. Each of the rounds is automatically loaded into a chamber in the firearm and ejected therefrom upon the firing thereof. The magazine-latching mechanism and in particular the control for operating the magazine-latching mechanism, depending on the type of firearm, is typically located on the handle or the frame of the firearm. Similarly, the control for operating the firing safety mechanism of prior known firearms typically comprises a non-self-retracting lever, a button or the like which is located on the frame of the firearm. The controls of the magazine-latching mechanism and the safety mechanism of prior known firearms, being in exposed positions, are subject to being inadvertently dislodged. If these controls are dislodged inadvertently, serious safety-related problems can result. For example, if the control for the magazine-latching mechanism is dislodged inadvertently, it could result in a dislodging of the magazine itself from the firearm during a firing of the firearm. Similarly, if the control for the safety mechanism is dislodged inadvertently, the firearm could be placed in a condition for firing at a time when an operator of the firearm believes that the safety mechanism is on.\nIn addition to firearm safety related features, the appearance of a firearm is also very important. In the minds of many users, a firearm with clean lines and a minimum of externally protruding controls is pleasant to view and a pleasure to operate."}
{"text": "The present invention pertains to electronic circuitry and particularly to a specialized circuit useful in telephone communication.\nIt is well-known in the telephone communication arts to provide at a telephone sub-station a suitable dial or other device for interrupting a DC path on a repetitive basis so as to transmit dial pulses that affect switching systems at a central station. It happens that in certain types of telephone systems it is required that the bell, buzzer, or other audible means used to indicate an incoming call, be continuously in circuit with the dial means or other instrumentality that is connected across the lines.\nThe difficulty or problem that arises in a telephone subset context when the aforesaid conditions apply is that so-called \"dial tap\" is encountered, that is to say, the buzzer, bell or other audible means for signaling an incoming call necessarily operates in response to the dial pulses placed on the line as the dial means is being rotated in effecting its purposes. In other words, an annoying sound is heard by the user of the telephone instrument when he is trying to dial out on the line.\nAccordingly, it is a primary object of the present invention to overcome the problem or difficulty of \"dial tap\" as it is encountered in certain telephone systems.\nAnother object is to accomplish the successful overcoming of the aforestated problem in a very simple, low-cost way.\nA further object is to provide a simple electronic circuit in a package that can be interposed between the conventional telephone lines and the buzzer device.\nThe above and other objects are achieved and implemented by a primary feature of the present invention according to which an electronic circuit that operates as a discriminating means is provided for insertion between the so-called \"line\" terminals of a telephone instrument and the \"device\" terminals, which are the terminals directly connected to a buzzer, bell or other audible device. This discriminating means functions so as to block the passage of the repetitive dial pulses and prevent them from affecting the buzzer device, while permitting conventional or regular ringing current of appropriate frequency to pass to the buzzer device thereby to actuate such device and to make audible to the user the fact that there is an incoming call.\nThe discriminating means operates by reason of a suitable time delay arrangement which is effective, in the case of the extended time period of the ringing signal, to actuate a latching device, such as a silicon controlled rectifier or similar device. On the other hand, since the dial pulse signals that would ordinarily cause the annoying dial tap problem are of short duration, that is, on the order of less than 100 milliseconds, these will not be effective to cause actuation of the latching device and hence the path for transmission of such dial pulses will be completely blocked. Moreover, means are provided that prevent build-up of the succession or series of dial pulses from affecting the latching device and causing spurious operation.\nOther and further objects, advantages and features of the present invention will be understood by reference to the following specification in conjunction with the annexed drawing, wherein like parts have been given like numbers."}
{"text": "1. Field of the Invention\nThe invention relates to a vehicle, and more particularly to an energy-saving vehicle.\n2. Description of the Related Art\nFor vehicles, upward resistance occurs during high-speed travelling, and automakers try to overcome the upward resistance by increasing weight of the vehicles. However, increase of weight also results in high energy consumption. Moreover, automakers use ethanol driven vehicles, hydrogen vehicles, air-powered vehicles and so on to replace conventional vehicles whereby saving energy. But an un-neglectable problem with them is that energy consumption thereof is still high after various transmission devices are used."}
{"text": "Over the past several years, internet-enabled smart phones, tablets and gaming consoles have become essential personal accessories, connecting users to friends, work, leisure activities and entertainment. Users now have more choices and expect to have access to content, data and communications at any time, in any place. As more users utilize these services, telecommunication networks must meet these increases in user demand, and support the array of new services and provide fast, reliable communications. Therefore, improved methods and solutions for dynamically allocating underutilized telecommunication resources (e.g., RF spectrum, etc.) of a first telecommunication network for access and use by wireless devices that subscribe to other networks will be beneficial to the telecommunication networks, service providers, and to the consumers of telecommunication services."}
{"text": "Aluminum is used as a material of wiring in various electronic devices such as a solar cell, a semiconductor device, and an electronic display device. Heretofore, it has generally been practiced to form such an aluminum film by a vacuum process such as a sputtering, vacuum deposition or chemical vapor deposition (CVD) method and to form a pattern by a photoetching method. Since this method requires a large-scale vacuum deposition device, not only is it disadvantageous in terms of energy consumption, but it is also difficult to uniformly form aluminum wiring on a substrate having a large surface area, resulting in poor yield and one of the causes of high costs.\nIn contrast, there has recently been proposed the technology in which a solution composition containing a complex of an amine compound and an aluminum hydride, and an organic solvent is applied on a substrate and then an aluminum film is formed by heating the solution composition (specification of Japanese Patent No. 3,817,201). This technology is an excellent technology which can form a satisfactory quality aluminum film by a simple method without requiring a massive and heavy apparatus. However, since this technology requires a pretreatment in which an organometallic compound of Ti, Pd or Al is applied on a substrate and then subjected to a heat treatment before application of the solution composition, the technology has such a drawback that a film made of an oxide of the metal is interposed in the space between the thus formed aluminum film and the substrate, thus failing to ensure electrical bonding between the substrate and the aluminum film."}
{"text": "The present invention relates to speech coders and, more particularly, to speech coders for high quality coding of speech signals at low bit rates.\nA speech coder is used together with a speech decoder such that the speech is coded by the coder and decoded in the speech decoder. A well known method of high efficiency speech coding is CELP (Code Excited Linear Prediction coding) as disclosed in, for instance, M. Schroeder, B. Atal et al, xe2x80x9cCode-Exited Linear Prediction: High Quality Speech at very low bit ratesxe2x80x9d, IEEE Proc. ICASSP-85, 1985, pp. 937-940 (Reference 1) and Kleijn et al, xe2x80x9cImproved Speech Quality and Efficient Vector Quantizatin in SELPxe2x80x9d, IEEE Proc. ICASSP-88, 1988, pp. 155-158 (Reference 1). In this method, on the transmission side, a spectral parameter, representing a spectral energy distribution of a speech signal, is extracted from the speech signal for each frame (of 20 ms, for instance) by using linear prediction (LPC) analysis. Also, the frame is further divided into a plurality of sub-frames (of 5 ms, for instance), and parameters (i.e., delay parameter corresponding to pitch period and gain parameter) are extracted for each sub-frame on the basis of the past excitation signals. Then, pitch prediction of a pertinent sub-frame speech signal is executed by using an adaptive codebook. For an error signal which is obtained as a result of the pitch prediction, an optimum excitation codevector is selected from an excitation codebook (or vector quantization codebook) constituted by a predetermined kind of noise signal, whereby an optimal gain is calculated for excitation signal quantization. The optimal excitation codevector is selected so as to minimize the error power between a signal synthesized from the selected noise signal and the error signal noted above. Index and gain, representing the kind of the selected codevector, are transmitted together with the spectral parameter and adaptive codebook parameter to a multiplexer. Description of the receiving side is omitted.\nIn the above prior art speech coder, enormous computational effort is required for the selection of the optimal excitation codevector from the excitation codebook. This is so because in the method according to References 1 and 2, described above the excitation codevector selection is executed by repeatedly performing for each codevector, filtering or convolution a number of times corresponding to the number of the codevectors stored in the codebook. For example, where the bit number of the codebook is B and the dimension number is N, denoting the filter or impulse response length in the filtering or convolution by K, a computational effort of Nxc3x97Kxc3x972Bxc3x978,000/N per second is required. By way of example, assuming B=10, N=40 and K=10, it is necessary to execute the computation 81,920,000 times per second. The computational effort is thus enormous and economically unfeasible.\nHeretofore, various methods of reducing the computational effort necessary for the excitation codebook retrieval have been proposed. For example, an ACELP (Algebraic Code-Excited Linear Prediction) system has been proposed. The system is specifically treated in C. Laflamme et al, xe2x80x9c16 kbps Wideband Speech Coding Technique based on Algebraic CELPxe2x80x9d, IEEE Proc. ICASSP-91, 1991, pp. 13-16 (Reference 3). According to Reference 3, the excitation signal is expressed with a plurality of pulses, and transmitted with the position of each pulse represented with a predetermined number of bits. The amplitude of each pulse is limited to +1.0 or xe2x88x921.0, and it is thus possible to greatly reduce the computational effort of the pulse retrieval.\nThe method according to Reference 3, however, has a problem that the speech quality is insufficient, although great reduction of computational effort is attainable. The problem stems from the fact that each pulse can take only either positive or negative polarity and that its absolute amplitude is always 1.0 irrespective of its position. This results in very coarse amplitude quantization, thus deteriorating the speech quality.\nAn object of the present invention is to provide a speech coder capable of preventing speech quality deterioration with relatively less computational effort even where the bit rate is low.\nAccording to the present invention, there is provided a speech coder comprising a spectral parameter calculation unit for obtaining a spectral parameter i.e. spectral energy distribution, from an input speech signal and quantizing the obtained spectral parameter, an excitation quantization on unit for quantizing excitation signal of the speech signal by using the spectral parameter and outputting the quantized excitation signal, the excitation being constituted by a plurality of non-zero pulses. The speech coder further comprising a codebook for simultaneously quantizing one of two, i.e., amplitude and position, paramters of the non-zero pulses, the excitation quantization unit having a function of quantizing the non-zero pulses by obtaining the other parameter by retrieval of the codebook.\nThe excitation quantization unit has at least one specific pulse position for taking a pulse thereat that is, a range within which at least one pulse position is to be determined.\nThe excitation quantization unit preliminarily selects a plurality of codevectors from the codebook and executes the quantization by obtaining the other parameter by retrieval of the preliminarily selected codevectors.\nAccording to another embodiment of the present invention, there is provided a speech coder comprising a spectral parameter calculation unit for obtaining a spectral parameter from an input speech signal for every frame and quantizing the obtained spectral parameter, and an excitation quantization unit for quantizing an excitation signal of the speech signal by using the spectral parameter and outputting the quantized excitation signal. The excitation signal is constituted by a plurality of non-zero pulses. The speech coder further comprising a codebook for simultaneously quantizing the amplitude of the non-zero pulses and a mode judgement circuit for executing mode judgement by extracting a feature quantity from the speech signal. The excitation quantization unit providing, when a predetermined mode is determined as a result of the mode judgement in the mode judgement circuit, functions of calculating positions of non-zero pulses for a plurality of sets, executing retrieval of the codebook with respect to the pulse positions in the plurality of sets and executing excitation signal quantization by selecting a combination of a codevector and a pulse position, at which a predetermined equation has a maximum or a minimum value.\nAccording to another embodiment of the present invention, there is provided a speech coder comprising a spectral parameter calculation unit for obtaining a spectral parameter from an input speech signal for every frame and quantizing the obtained spectral parameter, and an excitation quantization unit for quantizing an excitation signal of the speech signal by using the spectral parameter and outputting the quantized excitation signal. The excitation signal being constituted by a plurality of non-zero pulses. The speech coder further comprising a codebook for simultaneously quantizing the amplitude of the non-zero pulses and a mode judgement circuit for making a mode judgement by extracting a feature quantity from the speech signal. The excitation quantization unit providing, when a predetermined mode is determined as a result of the mode judgement in the mode judgement circuit, functions of calculating positions of non-zero pulses for at least one set, executing retrieval of the codebook with respect to pulse positions of a set having a pulse position, at which a predetermined equation has a maximum or a minimum value, and effecting excitation signal quantization by selecting the optimal combination satisfactory pulse position set and codevector. When a different predetermined mode is determined, the excitation quantization unit provides functions of representing the excitation in the form of linear coupling of a plurality of pulses and excitation codevectors selected from the excitation codebook, and executing excitation signal quantization by making retrieval of the pulses and the excitation codevectors.\nAccording to a further embodiment of the present invention, there is provided a speech coder comprising a frame divider for dividing input speech signal into frames having a predetermined time length, a sub-frame divider for dividing each frame speech signal into sub-frames having a time length shorter than the frame, a spectral parameter calculator which receives a series of frame speech signals outputted from the frame divider, truncates the speech signal by using a window longer than the sub-frame time and does spectral parameter calculation up to a predetermined degree. The speech coder further comprises a spectral parameter quantizer which vector quantizes a LSP parameter of a predetermined sub-frame, calculated in the spectral parameter calculator, by using a linear spectrum pair parameter codebook, a perceptual weight multiplier which receive line prediction coefficients of a plurality of sub-frames, calculated in the spectral parameter calculator, and does perceptual weight multiplication of each sub-frame speech signal to output a perceptual weight multiplied signal. The speech coder also includes a response signal calculator which receives, for each sub-frame, linear prediction coefficients of a plurality of sub-frames calculated in the spectral parameter calculator and linear prediction coefficients restored in the spectral parameter quantizer, calculates a responses signal for one sub-frame and outputs the calculated response signal to a subtracter. The speech coder further includes an impulse response calculator which receives the restored linear prediction coefficients from the spectral parameter quantizer and calculates an impulse response of a perceptual weight multiply filter for a predetermined number of points. An adaptive codebook circuit receives past excitation signal fed back from the output side, the output signal of a subtracter and the perceptual weight multiply filter impulse response, and obtains a delay corresponding to the pitch and outputs index representing the obtained delay. An excitation quantizer does calculation and quantization of one of parameters of a plurality of non-zero pulses constituting an excitation by using an amplitude codebook for collectively quantizing other parameter, i.e., amplitude parameter, of excitation pulses. A gain quantizer reads out gain codevectors from a gain codebook, selects a gain codevector from amplitude codevector/pulse position data and outputs index representing the selected gain codevector to a multiplexer. A weight signal calculator receives the output of the gain quantizer, reads out a codevector corresponding to the index and obtains a drive excitation signal.\nOther objects and features will be clarified from the following description with reference to attached drawings."}
{"text": "The circadian clock links our daily cycles of sleep and activity to the external environment. Deregulation of the clock is implicated in a number of human disorders, including depression, seasonal affective disorder, and metabolic disorders. Circadian rhythms are controlled in mammals by the master clock located in the suprachiasmatic nucleus of the hypothalamus (Antle and Silver, Trends Neurosci 28: 145-151). At the cellular level, the molecular events behind clock cycling are described by the regular increase and decrease in mRNAs and proteins that define feedback loops, resulting in approximately 24 hour cycles. The suprachiasmatic nucleus is primarily regulated, or entrained, directly by light via the retinohypothalamic tract. The cycling outputs of the suprachiasmatic nucleus, not fully identified, regulate multiple downstream rhythms, such as those in sleep and awakening, body temperature, and hormone secretion (Ko and Takahashi, Hum Mol Gen 15: R271-R277.). Furthermore, diseases such as depression, seasonal affective disorder, and metabolic disorders, may have a circadian origin (Barnard and Nolan, PLoS Genet. 2008 May; 4(5): e1000040).\nPhosphorylation of circadian clock proteins is an essential element in controlling the cyclical rhythm of the clock. CK1ε and CK1δ are closely related Ser-Thr protein kinases that serve as key clock regulators as demonstrated by mammalian mutations in each that dramatically alter the circadian period. (Lowrey et al., Science 288: 483-492). Therefore, inhibitors of CK1δ/ε have utility in treating circadian disorders. Thus it is an object of this invention to provide compounds of Formula I that are inhibitors of CK1δ or CK1ε. This object and other objects of this invention become apparent from the detailed discussion of the invention that follows."}
{"text": "1. Field\nThe present general inventive concept relates to an image forming apparatus supporting Wi-Fi direct, and more particularly, to an apparatus and method of activating Wi-Fi direct in an image forming apparatus supporting Wi-Fi direct.\n2. Description of the Related Art\nRecently, peer to peer (P2P) communication technology, by which wireless devices may be directly connected to each other without an additional wireless connecting apparatus, have been generalized and widely used. For example, Bluetooth technology makes the P2P communication possible. Although there are limitations in terms of transmission speed and transmission range of Bluetooth, new versions of Bluetooth are being developed to compensate for the limitations.\nIn addition, Wi-Fi, which is a wireless local area network (WLAN) standard based on IEEE 802.11 regulated by the Wi-Fi alliance, is basically a technology using ultra-high speed Internet by accessing access points (APs) connected to an infrastructured network; however, the Wi-Fi may serve the P2P communication by using an ad-hoc function. However, when the ad-hoc function is used, security is weakened, a transmission speed is lowered, and a setting method is not easily performed. Therefore, the Wi-Fi alliance has suggested a Wi-Fi Direct technology that makes the P2P communication possible. The Wi-Fi Direct allows P2P connection between wireless devices without using the AP, supports a transmission speed of a maximum of 250 Mbps, and performs security settings by using Wi-Fi protected access 2 (WPA2), in order to address problems of the ad-hoc function. In addition, the Wi-Fi Direct supports a transmission range of a maximum of 200 m, and thus, is considered as a substitute for the P2P communication.\nAs described above, with the appearance of the Wi-Fi Direct, it is considered that utilization of the P2P communication is increased more. In addition, the P2P communication technology may be also applied to image forming apparatuses such as printers, scanners, facsimiles, and multi-function printers. Therefore, technologies for safely and conveniently using an image forming apparatus supporting the P2P connection are necessary."}
{"text": "The present invention relates to lighting devices or illuminating lamps which have special advantage for personal use. That is, the devices may be used in association with the person, such as strapped to a wrist, or hung on a belt, or hung in a convenient place on a bicycle, baby carriage, or the like.\nNumerous and varied illuminating devices have been contrived and proposed in the prior art. Some of those devices have limited utility. Other of the prior devices may present undesirable complexities, or manufacturing difficulties, or other undesirable features."}
{"text": "Network nodes forward data. Network nodes may take form in one or more routers, one or more bridges, one or more switches, one or more servers, or any other suitable communications processing device. The data is commonly formatted as packets and forwarded using forwarding tables. A packet is a formatted unit of data that typically contains control information and payload data. Control information may include: information that identifies sources and destinations, such as addresses, error detection codes like checksums, sequencing information, etc. Control information is typically found in packet headers and trailers. Payload data is typically located between the packet headers and trailers.\nForwarding packets involves various processes that, while simple in concept, can be complex. The processes involved in forwarding packets vary, depending on the type of forwarding method used. Multicast is the preferred method of data forwarding for many networks. One reason for this is that multicast is a bandwidth-conserving technology that reduces traffic by simultaneously delivering data to multiple receivers. However, some network environments are not well suited to support multicast. Doing so in such environments often involves discovering and maintaining significant amounts of control, or state, information. Setting up and maintaining this control information has a tendency to become complex and costly in terms of computing resources, and can become a major limiting factor in overall network performance. Another issue with multicast is that due to packet delivery mechanisms used, packets are sometimes forwarded to locations where the packets were not desired. This unnecessary delivery of packets represents an unwelcome burden on network performance. Overcoming this burden by traditional means involves generation and maintenance of even more control information."}
{"text": "Many diagnostic systems require the capturing and analysis of a peak value of an input waveform. In an analog environment, a conventional Sample and Hold circuit (S/H) is commonly used to determine the peak value. However, S/H circuits typically have associated minimum input voltage requirements. In many diagnostic applications, where the input signal has a peak value below the minimum required input voltage, inaccuracies result due to distortion of the S/H circuit. Specifically, below the minimum input voltage, the input signal tends to become distorted by the noise of the S/H circuit, that might be comparable in magnitude to the peak, as shown by FIG. 1. As a result, the S/H circuit is incapable of reliably detecting peak values. Moreover, for waveforms having a first portion having a first rate of change and a second portion having a relatively faster rate of change, (or vice versa) as shown by FIG. 1, the S/H circuit may be incapable of capturing the peak values due to the second portion whose rate of change cannot be detected by the current.\nIn an analog environment, one solution to overcome the problem is to amplify the signal at a predetermined gain so that the S/H circuit functions reliably. Alternatively, a phase shift circuit may be implemented to detect the falling edge of the input signal, VIN. However, as phase shift circuits are plagued with many of the same problems as S/H circuits, they often tend to fail to function reliably. The remaining corrective measure is amplification of the input signal. In situations where power consumption must be kept to a minimum, such as battery powered diagnostic equipment, amplifying input signals consumes unnecessary amounts of power.\nOne way to address this problem is by digitally signal sampling the waveform with the use of an analog-to-digital (A/D) converter in order to capture the peak value. As illustrated by FIG. 2, the waveform falls abruptly at a rate faster than the sampling rate. In this regard, a peak value is calculated from the rising edge only. However, as a sample reading does not always capture the peak value, results will tend to be inaccurate. A high sampling rate may be used to overcome this problem, but is undesirable in many diagnostic system applications, as fast A/D converters cost more and require more power to operate. Especially for battery powered diagnostic systems in the field, fast A/D converters are not practical.\nAccordingly, a need arises to reliably detect peak voltages of a waveform, especially, an input signal level having high rates of change and/or are low in magnitude."}
{"text": "Various industrial manufacturing processes often require the use of gasses and fluids which are controlled by systems made up of valves, regulators, pressure transducers, mass flow controllers and the like. These components are typically connected together by the use of welded tubing and compression fittings and mounted on a vertical panel. These type of connections may be undesirable in some applications because they add additional time and cost for welding operations, unnecessary space between components and make it difficult to replace a component located between other components. Further, these systems are typically custom designed and manufactured which make the manufacturing costs and procurement of replacement parts quite expensive.\nNew modular fluid systems have been recently introduced into the semiconductor industry in order to overcome these type of problems. Typical components of these systems such as valves, pressure regulators and other typical fluid components have been reconfigured so that their inlet and outlet ports are co-located in a coplanar configuration. Further, the attachment flow component flange has a standard size and shape in order to permit interchangeability of surface mount components. However, these fluid systems have the disadvantage of being very expensive because they are machined from high purity metal stock. These systems further require the use of metal seals, which are very expensive. Thus it is desired to provide an inexpensive modular manifold system for use for example, in the analytical process industry.\nOther features and advantages of the invention will become apparent from the following detailed description, with reference to the accompanying drawing and claims, which form a part of the specification."}
{"text": "Relatively recently, researchers observed that double stranded RNA (“dsRNA”) could be used to inhibit protein expression. This ability to silence a gene has broad potential for treating human diseases, and many researchers and commercial entities are currently investing considerable resources in developing therapies based on this technology.\nDouble stranded RNA induced gene silencing can occur on at least three different levels: (i) transcription inactivation, which refers to RNA guided DNA or histone methylation; (ii) siRNA induced mRNA degradation; and (iii) mRNA induced transcriptional attenuation.\nIt is generally considered that the major mechanism of RNA induced silencing (RNA interference, or RNAi) in mammalian cells is mRNA degradation. Initial attempts to use RNAi in mammalian cells focused on the use of long strands of dsRNA. However, these attempts to induce RNAi met with limited success, due in part to the induction of the interferon response, which results in a general, as opposed to a target-specific, inhibition of protein synthesis. Thus, long dsRNA is not a viable option for RNAi in mammalian systems.\nMore recently it has been shown that when short (18-30 bp) RNA duplexes are introduced into mammalian cells in culture, sequence-specific inhibition of target mRNA can be realized without inducing an interferon response. Certain of these short dsRNAs, referred to as small inhibitory RNAs (“siRNAs”), can act catalytically at sub-molar concentrations to cleave greater than 95% of the target mRNA in the cell. A description of the mechanisms for siRNA activity, as well as some of its applications are described in Provost et al. (2002) Ribonuclease Activity and RNA Binding of Recombinant Human Dicer, EMBO J. 21(21): 5864-5874; Tabara et al. (2002) The dsRNA Binding Protein RDE-4 Interacts with RDE-1, DCR-1 and a DexH-box Helicase to Direct RNAi in C. elegans, Cell 109(7):861-71; Ketting et al. (2002) Dicer Functions in RNA Interference and in Synthesis of Small RNA Involved in Developmental Timing in C. elegans; Martinez et al., Single-Stranded Antisense siRNAs Guide Target RNA Cleavage in RNAi, Cell 110(5):563; Hutvagner & Zamore (2002) A microRNA in a multiple-turnover RNAi enzyme complex, Science 297:2056.\nFrom a mechanistic perspective, introduction of long double stranded RNA into plants and invertebrate cells is broken down into siRNA by a Type III endonuclease known as Dicer. Sharp, RNA interference—2001, Genes Dev. 2001, 15:485. Dicer, a ribonuclease-III-like enzyme, processes the dsRNA into 19-23 base pair short interfering RNAs with characteristic two base 3′ overhangs. Bernstein, Caudy, Hammond, & Hannon (2001) Role for a bidentate ribonuclease in the initiation step of RNA interference, Nature 409:363. The siRNAs are then incorporated into an RNA-induced silencing complex (RISC) where one or more helicases unwind the siRNA duplex, enabling the complementary antisense strand to guide target recognition. Nykanen, Haley, & Zamore (2001) ATP requirements and small interfering RNA structure in the RNA interference pathway, Cell 107:309. Upon binding to the appropriate target mRNA, one or more endonucleases within the RISC cleaves the target to induce silencing. Elbashir, Lendeckel, & Tuschl (2001) RNA interference is mediated by 21- and 22-nucleotide RNAs, Genes Dev. 15:188, FIG. 1.\nThe interference effect can be long lasting and may be detectable after many cell divisions. Moreover, RNAi exhibits sequence specificity. Kisielow, M. et al. (2002) Isoform-specific knockdown and expression of adaptor protein ShcA using small interfering RNA, J. Biochem. 363:1-5. Thus, the RNAi machinery can specifically knock down one type of transcript, while not affecting closely related mRNA. These properties make siRNA a potentially valuable tool for inhibiting gene expression and studying gene function and drug target validation. Moreover, siRNAs are potentially useful as therapeutic agents against: (1) diseases that are caused by over-expression or misexpression of genes; and (2) diseases brought about by expression of genes that contain mutations.\nSuccessful siRNA-dependent gene silencing depends on a number of factors. One of the most contentious issues in RNAi is the question of the necessity of siRNA design, i.e., considering the sequence of the siRNA used. Early work in C. elegans and plants circumvented the issue of design by introducing long dsRNA (see, for instance, Fire, A. et al. (1998) Nature 391:806-811). In this primitive organism, long dsRNA molecules are cleaved into siRNA by Dicer, thus generating a diverse population of duplexes that can potentially cover the entire transcript. While some fraction of these molecules are non-functional (i.e., induce little or no silencing) one or more have the potential to be highly functional, thereby silencing the gene of interest and alleviating the need for siRNA design. Unfortunately, due to the interferon response, this same approach is unavailable for mammalian systems. While this effect can be circumvented by bypassing the Dicer cleavage step and directly introducing siRNA, this tactic carries with it the risk that the chosen siRNA sequence may be non-functional or semi-functional.\nA number of researches have expressed the view that siRNA design is not a crucial element of RNAi. On the other hand, others in the field have begun to explore the possibility that RNAi can be made more efficient by paying attention to the design of the siRNA. Unfortunately, none of the reported methods have provided a satisfactory scheme for reliably selecting siRNA with acceptable levels of functionality. Accordingly, there is a need to develop rational criteria by which to select siRNA with an acceptable level of functionality, and to identify siRNA that have this improved level of functionality, as well as to identify siRNAs that are hyperfunctional."}
{"text": "Data stored on one or more computer systems may be useful for a user. However, the stored data may be too extensive for a user to find the data by direct examination. Additionally, some parts of a data repository may contain information that is not accessible to the user. In many cases, in order to allow the user useful access to the data, a search mechanism is provided. The search mechanism allows a user to issue a search request (also termed a search query). The query is executed and the results are returned for the user.\nFor example, a web-based search engine is a search mechanism that is used to provide search access to information via a web-based search. The information may be found in a specific data repository, such as a database or other data collection. The information may also be an agglomeration of data found in a number of different data repositories. Such a search engine may provide search access to information available from different information providers over a network, such as the Internet.\nIn a typical usage of a web search engine, the user enters a query, which is a set of search terms related to what the user is looking for. The query is transmitted to the search engine, which attempts to locate “hits”—i.e., content that is available on the Internet and that relates to the terms contained in the query. Generally, the search engine either has a database of web pages that are known to exist, or communicates with external “providers” who maintain such databases; the query is “scored” against items in these databases to identify the web pages that best match the query. A list of results is then generated, and these results are returned to the user's computer for display by the user's web browser.\nTypically, the databases contain information such as: the Uniform Resource Locators (URLs) of web pages, the titles of the pages, descriptions of the pages, and possibly other textual information about the web pages. The user then reads the results and attempts to determine, based on the text contained in the results, whether the results correspond to what the user is looking for. Users may then attempt to retrieve the entire page correlating to a search result. In other contexts, search engines present results summarizing the pieces of data that may possibly be useful for a user.\nThe utility of the search engine is correlated directly to the quality of the results provided. In the best case, the results are presented to the user in order of utility to the user on the result page. Because the quality of the results is subjective, the user's satisfaction must be determined in order to determine whether the quality of the results were satisfactory.\nIn the prior art, quality of individual web pages has been measured by obtaining explicit feedback from a user. At least one prior art web browser has attempted to obtain such explicit feedback from a user. This browser is described in a paper entitled “Inferring User Interest” by Mark Claypool, David Brown, Phong Le, Makoto Waseda in IEEE Internet Computing 5(6): 32-39 (2001). In that browser, different pages are displayed by the browser. Whenever the page being displayed by the browser is changed, a user evaluation of the page is requested from the user. User evaluations for a given page are collected, to determine whether users find that page valuable. In this browser, some implicit feedback is also maintained regarding each page, including data regarding the time spent on the page, mouse movements, mouse clicks, and scrolling time.\nWhile this technique does gather user feedback, it has limited utility in situations in which users may have different needs for a page. For example, a user who is looking for information about books written by Douglas Adams may evaluate a page on his book The Hitchhiker's Guide to the Galaxy and give a high score for utility. However, another user who is looking for information on books about traveling cheaply may evaluate the same page and give it a low score. Thus the technique described will have limited utility in the wide variety of situations in which different users may have different needs, or even where a single user may have different needs for information at different times. In other words, the usefulness of this technique is limited because evaluation of each page is completely independent of the context in which the user arrived at the page.\nThus, this technique is not useful for evaluating the quality of a search engine. In general, this technique is not useful for evaluations that are context-based, but only for evaluating the quality of individual data items, independent of the context in which a user arrived at the data items.\nAnother drawback of the prior art is that users may respond to explicit requests for feedback by ignoring them or providing inaccurate feedback in order to speed the searches the users are performing. This leads to possibly unreliable data in addition to user dissatisfaction with the search tool.\nAdditionally, the prior art web browsers do not request feedback on non-standard search results. Many web searches now provide non-standard search results. For example, when searching web pages, for example, for a restaurant of a specified name in a specified town, some web searches simultaneously do a phone book search for the phone number of the restaurant. This result is displayed for the user. This is a non-standard result because it is not the type of result that the user was ostensibly requesting. However, it may be useful to judge the satisfaction of the user with such non-standard results.\nAn additional problem in judging satisfaction with search results is that they may be used by a user without the user selecting the result (standard or non-standard) and without the user otherwise indicating through other implicit feedback that the user has found a satisfactory result. In the above example, a user who is only looking for a web page regarding the restaurant to find the phone number of the restaurant may find that phone number in a non-standard result and use that information without selecting any result, standard or non-standard.\nAnother drawback of the prior art is that feedback from each user is considered without reference to the user making the feedback. A user with a different experience level, language, purpose, or technical ability may have different responses to the search mechanism, however there is no way to include such differences in considering user satisfaction according to the prior art systems.\nThus, there is a need for a system and method to overcome these deficits in the prior art. The present invention addresses the aforementioned needs and solves them with additional advantages as expressed herein."}
{"text": "The discussion below is merely provided for general background information and is not intended to be used as an aid in determining the scope of the claimed subject matter.\nAspects of the present invention relate to equipment and methods used during laboratory testing of bodies moving over surfaces such as but not limited to all types of vehicles traveling on a roadway.\nReferring to vehicles such as automobiles by way of example, the use of simulated roadways has become important in the development of vehicle designs. These roadway simulators generally comprise one or more individual endless belts that are mounted below the vehicle, engaging one or more wheels, and/or positioned below other body panels of the vehicle. The belts are driven to simulate driving the vehicle over a road. Commonly, the roadway simulators are disposed in wind tunnels so as to measure aerodynamic characteristics. Generally, roadway simulators improve the realism of flow characteristics leading to improved measurement of aerodynamic characteristics.\nThe vehicle is typically restrained with respect to the simulated roadways through a linkage that restrains the vehicle from movement in longitudinal direction (fore and aft), in lateral direction (side-to-side) and in yaw (steering motion). Commonly, the vehicle restraints are passive restraints which do not affect the vehicle's dynamics, but which can be used for measuring loads if desired in the different directions of restraint."}
{"text": "The goal of hybrid development is to combine, in a single hybrid, various desirable traits. For field crops, these traits may include resistance to diseases and insects, resistance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination, stand establishment, growth rate, maturity, and plant and ear height is important. Traditional plant breeding is an important tool in developing new and improved commercial crops."}
{"text": "It is known in the art to shrink form containers such as drinking cups from preformed tubular lengths of circumferentially oriented thermoplastic material such as foamed polystyrene.\nOne particularly desirable method of initially forming a tubular length of such circumferentially oriented material is to provide rectangular-preprinted blanks and wrap these blanks around a mandrel whereon a heat sealed seam is effected longitudinally along the circumference of the formed tubular length. The use of rectangular blanks facilitates pre-printing of patterns, designs, logos, etc., on the blanks such that the ultimate tubular lengths and containers formed therefrom will bear the ultimately desired indicia.\nA further advantage of the rectangular blank is that it may be cut from an extruded sheet of thermoplastic or thermoplastic foam which is stretched longitudinally, i.e., in the most logical, natural and facile direction of stretch after extrusion, namely, the machine direction, to achieve the necessary circumferential orientation in a tubular length or cylinder formed from the rectangular blank.\nPrevious efforts to handle these rectangular blanks and form them into cylinders, however, have required relatively elaborate systems of transfer rollers, turrets with multiple mandrels thereon and vacuum systems to properly index leading and/or trailing edges of the rectangular blanks on the transfer rollers or mandrels."}
{"text": "1. Field of the Invention\nThis invention relates to floor care appliances and, more specifically, relates to nozzle configurations for nozzles utilized in such floor care appliances.\n2. Description of the Prior Art\nAlthough attempts at providing constant air velocity and pressure nozzles are broadly old, no one heretofore is known to have utilized the constant pressure, constant velocity structure to also actually direct the air and entrained dirt flow towards the suction connection. Thus, these two concepts are combined in a nozzle so that the structure utilized for the combined functions is one and the same and so that the combined cleaning effect afforded by this structure is additive as to each and not subtractive as to either.\nAccordingly, it would seem advantageous to provide such a nozzle structure having the advantages of both substantially uniform air suction and dirt pick up, heightened by a directed air path configuration.\nIt would be still a further advantage to utilize a single composite nozzle structure in which the unitary configuration of it would provide both of these advantages.\nIt would be still further advantageous to provide a nozzle with an angled ledge which would yield in conjunction with the agitator, a directed a effect to the air and entrained dirt to move it towards the suction tube connection for the nozzle.\nIt would, additionally, be advantageous to provide an expanding groove to increase the cross sectional area of the nozzle, as it approached its suction tube, to provide, as much as possible, for a uniform velocity and pressure of suction air across the nozzle mouth to promote uniform cleaning."}
{"text": "High quality tape decks such as the ones used in sound studios require frequent adjustments in order to optimize performance. Two of the most common adjustments are adjusting the azimuth of the recording and playback head and adjusting the recording and playback levels. Although these adjustments are necessary on a very frequent basis, preferably before each recording session in a sound studio, these adjustments have required expensive instruments and a relatively high degree of expertise in adjusting tape decks and operating laboratory instruments.\nIn order to adjust the playback voltage level, the conventional technique is to play a test tape in the tape player. The test tape will have numerous tones recorded on it at a magnetic flux density that should produce a nominal voltage level output. Most tape players have a nominal output voltage of +4 dbm (1.23 volts), but -8 dbm (0.3 volts) and +2 dbm (1 volt) are also common nominal record and play levels for professional recorders. In order to measure these relatively low voltage readings, expensive volt meters are required when performing this adjustment using a conventional technique.\nIn some professional tape recorders, the recording head may be switched into a playback mode and the test tape may be played back through the recording head. Thus, the record level may be adjusted in the same manner as the playback level. However, in machines that do not have such feature, the record level may be adjusted by using an oscillator to produce an input signal of a selected frequency and voltage level which is recorded by the recording head on a tape and simultaneously played back through the playback head. Assuming the playback level has previously been adjusted, the output from the playback head will indicate the recording level of the tape deck and the record level may be adjusted to an appropriate level by monitoring the playback level. In order to perform this latter adjustment technique, an accurate oscillator must be available to the operator in order to produce the test signal.\nThe conventional technique for adjusting the azimuth of a playback head is to play a test tape having a manaural recording on the outermost two tracks of the tape. If the azimuth of the head is not properly aligned, the playback output signals from these two outermost tracks will be out of phase. In order to adjust the azimuth of the playback head, the playback output signal from the two outermost tracks are typically monitored on a single oscilloscope using two probes. The azimuth of the playback head is then adjusted until the output signals are phase synchronized.\nThe conventional technique for adjusting the azimuth of a record head is to input a manaural signal into the tape recorder, record this test signal on the two outermost tracks of a tape and play back these two outermost tracks through the playback head. Assuming that the azimuth of a playback had been previously adjusted, any phase difference between the two playback signals will be caused by the maladjusted azimuth of the record head. Thus, the azimuth of the record head is adjusted until the output playback signals from the two outermost tracks of the tape are synchronized. Again, conventionally, a two probe oscilloscope was necessary in order to monitor the relative phases of the two playback signals from the outermost tracks of the tape.\nThe primary drawback with the conventional techniques described above is that expensive laboratory equipment is required in order to perform the various adjustments. Also, considerable expertise in operating the test equipment and the tape deck were required in order to perform the adjustments, and most lab instruments have a visual output so that the technician must constantly watch the instruments while he is attempting to adjust the tape deck. When the tape deck's adjustment mechanisms are difficult to access, it is inconvenient for the operator to constantly view his instruments while attempting to adjust the tape deck.\nThus, a need has arisen for a simple and inexpensive instrument for use in adjusting the record and playback levels of a tape deck and for adjusting the azimuth of the recording and playback heads of a tape deck."}
{"text": "Common pipe fittings for liquid (or steam) have internal threads on the inside surface of the pipe-receiving end of the fitting head, and external threads on the outside surface of the matching end of said pipe. Then, said threaded pipe end is screwed into said pipe fitting with the engagement of said internal thread of the said pipe and the external thread of said pipe. Since any thread has depth, the mechanical strength of said threaded end will be greatly weakened. Since the thread on the end of pipe fittings are typically coarsely finished, screwing a liquid (or steam) pipe into a fitting is met with resistance. A common wrench is not useful to tighten common pipe fittings. Due to the smooth and round pipe body. Therefore, a pipe wrench with specially made chops should be used for biting on the pipe body. In doing so, the pipe can be screwed into the fitting. After assembling, scars and deformations caused by said wrench chops remain on the pipe body.\nFurthermore, when the assembled pipe, used for a considerable time is disassembled, the process is generally more difficult than assembling because, both threaded portions of the pipe and of the fitting bind tightly together due to rust, deformation and gluing etc. Therefore, after disassembling, the scars and deformations caused by said wrench chops will be made more serious. At the same time, it is possible the whole pipe will be broken at the threaded portion.\nOvertime, different piping materials have been used such as: wrought iron (zinc plated), plastics (hard PVC), brass, brass with plastic coating, wrought iron with inner plastic coating, copper with outer nickel-plating and alloy. Even if the external surface of the pipes is plated or coated, they are still not suitable to be bitten by pipe wrench because none of the plating or coating layers can bear the terrible biting of the wrench chops and any layer will be destroyed after the biting of wrench chops. Even copper pipes or alloy pipes with hardnesses lower than the hardness of wrench chops, will exhibit ugly scars and deformations on the pipe surface after a pipe wrench is used on the pipes.\nAs a result, many fittings that do not match by threading have been developed. One such pipe is provided with an integral flange on the pipe end. This flange method is only suitable for pipes made of soft materials such as brass. Another such pipe utilizes a plastic pressing ring located on the external surface of the pipe end. This method is usually used on indoor drain systems. Though the above mentioned methods have eliminated threads from pipe ends, the connection between the pipe and the fitting is still achieved by threading. The difference lies in that the external thread is made on the external surface of the fitting, and the internal thread is made on the inner surface of a casing head which sleeved on the pipe. Although these methods can solve the problem of pipe surface being damaged by wrench chops, but they have limited applications since they are not suitable for wrought iron pipes, and because the construction of threaded casing heads is complex. Furthermore, the process for flanging and screwing are troublesome.\nAnother solution was to replace the internal thread of prior art fittings with an internal teethed lock ring so that pipe without end threads can be directly inserted through the internal teethed lock ring integrated on the fitting and can be kept in position. Once it is in position, it is impossible to disassemble the pipe. But in practice, pipe pre-assembly is absolutely necessary. If an assembled set can not be disassembled, it means there can be no pre-assembly and no repair work is permissible. Therefore, nobody uses this method nowadays."}
{"text": "The present invention relates to external cavity lasers.\nAn external cavity laser comprises a cavity that provides an optical path so that a laser beam, stimulated by a gain or active medium, can resonate. The optical spectrum of an external cavity laser is limited by the spectral region over which the gain medium yields optical gain. Depending on the geometry of the cavity, the laser radiation has certain allowed spectral modes, also referred to as longitudinal modes. The spectral modes correspond to different distinct wavelengths or frequencies of the light. The allowed radiation spectrum in a laser cavity thus is a discrete distribution of modal frequencies limited by the gain curve. Examples of external cavity lasers can be found in U.S. Pat. No. 5,331,651, EP-A-0 552 394 and in EP-A-1 202 409.\nU.S. Pat. No. 5,331,651 discloses a tunable Fabry Perot etalon to be used in an external cavity laser comprising a semiconductor laser chip and adjustable diffraction grating. The adjustable diffraction grating is used for a coarse measurement of the laser output wavelength and the etalon is used for the fine adjustment."}
{"text": "U.S. Pat. No. 488,255 proposes a system for transmitting a complex signal comprising a data transmission signal component and an AM broadcast signal component. The transmitters using the method of Kahn, in particular, that described by patent U.S. Pat. No. 488,255 for the transmission of an AM broadcast signal, comprise separate means of processing of the envelope and of the phase and a multiplier of the envelope signal and of the phase signal at the output.\nIn digital modulation, the phase channel is not as in AM a pure RF frequency but a phase modulated RF frequency. As the envelope signal and the phase signal are processed by different and independent means, the delay between these two signals, envelope and phase, must be accurately tailored to obtain an RF output signal of high quality.\nA loop allows automatic elimination and dynamic adjustment of the delay between the envelope signal and the phase signal.\nThe technique generally used in loops is to compare in the time domain the initial baseband signal with the demodulated transmitted RF signal. To do this, the delay between the two signals, initial and demodulated transmitted, is evaluated accurately. Then, the envelope and phase parts, in digital mode, of the two signals, initial and demodulated transmitted, are compared to determine the delay between envelope and phase."}
{"text": "1. Field of the Invention\nThis invention relates to the field of heat barriers or shields for the wall behind free standing heating appliances as well as the floor upon which such heating appliances may restingly be supported.\n2. Description of the Prior Art\nIt is well-known that energy has been in an abundant supply in this country for the past several decades. Accordingly, during this time, conservation of energy, at least in most modern day residential structures, has been of little or no interest. However, the recent spiraling cost of fuel and the recognized necessities for conserving energy are bringing about changes in the manner in which many homes are heated today. Among these changes is the reintroduction of the old fashioned potbellied stove concept. Of course, many modern modifications to the old type stove have been made, but the presence of the free standing heating appliance or space heater, usually using wood, coal or other materials as a source of fuel, is becoming more and more common throughout the country.\nIt is well-known that these type heating appliances develop a tremendous amount of heat usually necessitating that the wall adjacent thereto be protected in some form or other and that the floor also be protected from this heat. Therefore, there exists a need for a thremoshield which is, of course, noncombustible. Moreover, the thermoshield should contain no asbestos.\nThe concept of providing a thermoshield for this purpose was first disclosed many years ago. In fact, applicant is aware of the following U.S. Pat. Nos. pertaining to various shields of this nature: 485,363 granted to Adrianson in 1892; and 1,764,595 granted to Anderson in 1930. Other patents pertaining to similar structure of this nature, as known by the applicant, consist of: U.S. Pat. No. 801,690 granted to Ralls in 1905; U.S. Pat. No. 874,296 granted to Brown in 1907; U.S. Pat. No. 1,970,504 granted to Nordan in 1934; U.S. Pat. No. 2,482,376 granted to Serkland in 1949; U.S. Pat. No. 2,629,370 granted to Janski in 1953; and U.S. Pat. No. 4,148,302 granted to Patterson in 1979. It is significant to note that none of the above mentioned patents suggests or discloses applicant's device."}
{"text": "The following disclosure relates generally to integrated circuit (IC) devices, and to methods and apparatus for the design and fabrication of IC devices.\nAs technology progresses and IC devices include smaller dimensions and increased feature density. As technology nodes shrink, challenges are raised for example, mismatch of device performance becomes more critical. However, for relatively larger devices typically required in system-on-a-chip (SOC), analog, digital signal processing (DSP) and radio frequency (RF) applications, it may be difficult to provide adequate matching of large devices in an IC. These device mismatch issues may arise from the replacement or gate-last methodology employed to provide metal gate technology. The gate-last methodology typically requires additional chemical mechanical processing (CMP) steps. These CMP processes can create differences in gate heights due to loading effects. These loading effects are often exacerbated by the use of large devices (e.g., in concert with smaller features). For example, CMP dishing can result which may result in metal work-function shifts and mismatch degradation of the IC.\nThus, what is needed is an improved manner of circuit design, fabrication and implementation for gate last processes of ICs including differently sized devices."}
{"text": "One type of wet cell battery commonly used in automotive vehicles includes a pair of terminals extending from the top or side of the battery case. A threaded stud extends from a relatively large diameter shoulder portion which protrudes a short distance from the case. The prior art terminations have incorporated a large area flat contact surface which rests on the surface of the terminal shoulder when the threaded stud extends through a mounting hole in the flat contact surface of the termination. A nut is threaded onto the stud to press the contact surface tightly against the shoulder.\nUnless the contact surface at the bottom of the termination is perfectly flat and parallel to the top surface against which the nut is tightened, an excessive torque is required to prevent rocking of the termination on the terminal and to provide a low ohomic connection between the termination and the battery terminal. The use of too great a torque can strip the threads on the terminal stud or otherwise damage the battery.\nBattery cable terminations have included a plastic housing or cover molded over the metallic parts of the termination and such terminations have had associated with them several types of seals for sealing the termination from the battery acid fumes and from the elements. U.S. Pat. No. 4,049,335 describes a battery cable termination in which the molded cover sealably engages the terminal shoulder to prevent the battery acid fumes from creeping along the terminal and corroding the termination. U.S. Pat. No. 4,126,367 describes a battery cable termination which includes a sealing bead which contacts the battery case throughout an area surrounding the terminal. It would be desirable to provide an improved seal to the battery case and which also prevents the battery acid fumes from reaching the metal parts of the termination."}
{"text": "The rate of chronic and acute wound healing can be delayed or impaired by a number of factors (exogenous and endogenous) and a variety of medical conditions. Examples include infection, ulceration particularly through diabetes, circulation problems associated with vascular disease, malnutrition, stress, cancer radiotherapy and/or chemotherapy, compromise of the immune system or simply due to the normal aging process. At present there is a clear need for therapeutic and cosmetic approaches that promote wound healing processes.\nThe literature describes the PDT use of various photosensitisers in attempts to improve wound healing. For example, Photofrin (a porphyrin, Photofrin is a trade mark of Johnson & Johnson) showed delayed wound healing in rat skin flap reconstruction (Kubler et al Lasers in surgery and Medicine (1996), 18(4), 397-405. Parekh et al concluded in Lasers in surgery and Medicine (1999), 24(5), 375-381 that a benzoporphyrin and a phthalocyanine didn't alter wound healing. Lambrects et al concluded in Photochemical & Photobiological Sciences (2005), 4(7), 503-509 that a porphyrin delayed wound healing. Hamblin concluded in Journal of infectious diseases (2003), 187(11), 1717-25 that for a photosensitiser (Poly-lysine-chlorin e6) topical PDT did not accelerate wound healing, compared with that in untreated wounds. Belmont et al in the Laryngoscope (1999), 109(6), 886-90 using Photofrin concluded that PDT has been shown to delay wound healing. Heckenkamp et al in Journal of vascular surgery, (2000), 31(6), 1168-77 concludes that local photodynamic action of methylene blue exerts a retarding effect on intimal hyperplasia (one of the natural wound healing processes).\nCompounds, methods and compositions/medicaments have now been identified which address these problems and provide a basis for promoting the wound healing process."}
{"text": "1. Field\nEmbodiments relate to a semiconductor device.\n2. Description of the Related Art\nRecently, a highly integrated semiconductor device including a transistor has been considered. The transistor may include a metal gate structure."}
{"text": "The National Fire Protection Association (NFPA) has promulgated standards—see NFPA 1971 Standard on Protective Ensemble for Structural Firefighting (2000 Edition)—for protective clothing for firefighters fighting structural fires. The National Fire Protection Association has promulgated standards—see NFPA 1977 Standard on Protective Clothing and Equipment for Wildland Firefighting (1998 Edition)—for protective clothing for firefighters engaged in fighting wildland fires. Presently, the National Fire Protection Association is contemplating standards for protective clothing for firefighters engaging in what are known as technical rescues, such as extracting victims from automobile wrecks.\nAs contemplated by the aforenoted standards for protective clothing for firefighters fighting structural fires, a protective garment has an outer shell, which provides abrasion resistance and puncture resistance, a liner including a moisture barrier, and a thermal liner. Conventionally, as known heretofore, the liner including the moisture barrier and the thermal liner are sewn together or are bonded together, and the sewn-together or bonded-together liners are attached detachably to and within the outer shell by suitable fasteners, such as zippers, hook-and-loop fasteners, and snap fasteners. Thus, if the moisture barrier fails while the thermal liner remains useful, the liner including the moisture barrier cannot be easily replaced, apart from the thermal liner.\nAlternatively, as known heretofore, the liner including the moisture barrier is sewn to and within the outer shell and the thermal liner is attached detachably to the outer shell, within the liner including the moisture barrier. Thus, if the moisture barrier fails while the outer shell remains useful, the moisture barrier cannot be easily replaced, apart from the outer shell.\nAs contemplated by the aforenoted standards for protective clothing for firefighters fighting wildland fires, a protective garment has an outer shell, which provides abrasion resistance and puncture resistance. Although a liner including a moisture barrier and a thermal liner are not needed, a thermal liner may be optionally included, as for wearing under cold winter conditions. If a thermal liner is included, the thermal liner is attached detachably to and within the outer shell by suitable fasteners, such as zippers, hook-and-loop fasteners, or snap fasteners.\nThe aforenoted standards contemplated for protective garments for firefighters engaging in technical rescues are expected to provide that a protective garment has an outer shell, which provides abrasion resistance and puncture resistance, and a liner including a moisture barrier. Furthermore, the contemplated standards are expected to permit the liner including the moisture barrier to be detachably attached to and within the outer shell by suitable fasteners, such as zippers, hook-and-loop fasteners, or snap fasteners.\nProtective garments conforming to the aforenoted standards for protective clothing for firefighters fighting structural fires and protective garments conforming to the aforenoted standards for protective clothing for firefighters fighting wildland fires are available commercially from Morning Pride Manufacturing, L.L.C. of Dayton, Ohio, and from other sources."}
{"text": "Water may be entrained in petroleum-based liquid fuel as free water suspended in droplet or particulate form, in emulsion, or in solution. Dissolved water in fuel is common and generally presents little difficulties for operation of particulate filters for the fuel, as long as the water remains uniformly distributed in the fuel. The capacity of fuel to hold water in solution changes as a result of temperature changes, and in many cases, the finely-divided water particles will coalesce into larger particles in an undissolved state, which settle out because their density is greater than the fuel. Particulate fuel filters, frequently being formed from paper, degrade as the pores of the paper become clogged with the undissolved water, which can increase the pressure drop across the filter.\nIt is therefore important to remove the undissolved water particles suspended in the fuel before the fuel passes through any filter medium for removing solid particulate contaminants. The removal of water from the fuel system also prevents filter and fuel line freeze up.\nIt is also important that a water separator for a fuel supply system for a diesel engine pass solid particles of a larger size than the droplets of the water to be removed. Under some low temperature conditions when the suspended droplets of water are removed from the fuel by freezing, wax components of the diesel fuel coagulate and may form wax particles. It is important that the water separator pass such particles to prevent clogging of the separator and consequently, engine stalling.\nOne early type of water separator shown in U.S. Pat. No. 3,785,970, where a nylon sock had its open end sealed around the inlet passage to the separator so that the flow of fuel through the sock maintain the sock in an inflated condition. Using a nylon sock, there is an uniform minimal pressure drop across the entire porous surface of the sock. The pores of the sock are expandable to reach a state of equilibrium in which pores are of a size to create the pressure drop required to accommodate the established flow rate of fuel. The expandability of the pores also permits the passage of filterable particulate matter suspended in the fuel so that the water separator is not clogged by large wax particles of fuel under low-temperature operating conditions.\nWhile the water separator shown in the '970 patent has certain satisfactory qualities, the water separator has certain limitations as to structure and application, and was not particularly suited for a spin-on capability with a replaceable transparent collection bowl, to allow the separator to be easily attached to and detached from a manifold head, and to visually inspect the level of water in the separator element. Such spin-on types of elements with collection bowls have received wide-spread acceptance in the marketplace and consumers are generally comfortable and knowledgeable about their use and application.\nIt is also believed there is a constant demand in the industry for new and improved water separators, particularly water separators which incorporate few components, are easy to assemble, and which provide efficient separation of water from a water/fuel mixture."}
{"text": "The present invention relates to a cutting device for machines for cutting hides and the like. More specifically, the invention relates to a cutting device that is adapted to be used in machines for cutting hides, leather, materials in sheets or rolls, etc.\nAs is known, in the sector of machinery for digitally controlled cutting of materials that are extended flat, sheet-like, natural or synthetic, such as hide, leather, materials in sheets or rolls, etc., the cutting tool generally used is mounted in a tool head and is constituted by a blade that is adapted to oscillate vertically at high frequency and is moved in a substantially horizontal plane by way of digital control, on the material along the path of shapes to be cut which are defined or projected on the material.\nSo as to be deployed correctly on the material to be cut, in order to pass from the cutting of one shape to the next and furthermore in order to execute cutting paths with sharp corner edges, which require the extraction of the blade from the material, it is necessary to bring the oscillating blade vertically out of the material to be cut, in order to then move it conveniently by way of rapid movements, imposed by the digital controller, through the movement axes.\nTo such end, the cutting device or cutting head must be provided with a system of vertical movement.\nThe aim is to be able to raise the cutting blade as quickly as possible, especially when executing cutting paths with sharp corner edges (such as for example jagged sections of the outline of the piece), and this is because it enables a reduction in the times for executing the cutting operation.\nCurrently, the conventional cutting devices have a vertical rise speed of the blade, by way of movement of an oscillating spindle or of an entire tool head, which is limited by the inertia of the moved assembly, by the characteristics of guide systems along which the assembly is moved, and by the dynamic characteristics of the actuators that are adapted to achieve the above mentioned movement.\nTechniques are known of returning pneumatic pistons to the stroke limit position, which are based on a spring inserted in the cylinder which pushes the piston into abutment in the absence of a pneumatic supply to the latter.\nOne solution for the above mentioned cutting devices which is based on such technique thus involves the insertion of a spring propelling the oscillating piston toward the upper stroke limit position, once the pneumatic supply for oscillation of the piston that supports the cutting blade has been deactivated. Such configuration has however a number of drawbacks including the fact that such spring would produce a constant force on the oscillating piston, thus limiting its power during oscillation and penalizing the dynamic characteristics in terms of the ratio of stroke to oscillation frequency.\nFurthermore, such spring would require a high capability to withstand extremely high operating frequencies, which is difficult and costly to obtain."}
{"text": "The communications industry is rapidly changing to adjust to emerging technologies and ever increasing customer demand. This customer demand for new applications and increased performance of existing applications is driving communications network and system providers to employ networks and systems having greater speed and capacity (e.g., greater bandwidth). In trying to achieve these goals, a common approach taken by many communications providers is to use packet switching technology. Increasingly, public and private communications networks are being built and expanded using various packet technologies, such as Internet Protocol (IP). Note, nothing described or referenced in this document is admitted as prior art to this application unless explicitly so stated.\nMulticast Routing protocols like Protocol Independent Multicast-Sparse Mode (PIM-SM) and Bi-directional PIM achieve multicast forwarding by building a shared distribution tree rooted at Rendezvous Point (RP). In each of the routers (or switches), a single (*, G) entry is used to forward traffic down the shared tree. In Bi-directional, the same entry is used for forwarding traffic upstream towards the Rendezvous point, as well. In such shared tree based forwarding models, source-control is one of the biggest problems. Because traffic forwarding is achieved using a (*, G) entry, any unwanted malicious host can send multicast traffic at very high rates. This at least results in a waste of valuable network bandwidth, and lends itself to denial of service (DOS) attacks (e.g., attacks on a network by flooding it with useless traffic designed to consume the network resources).\nIn PIM-SM, data traffic from a source is encapsulated within a register message by the first hop router and sent via unicast packets to the Rendezvous Point (RP). Depending on the IP address of the source, the RP can decide to either forward the traffic down the shared tree or ignore the traffic. In either case, the RP sends back a register-stop message instructing the first hop router to stop sending register messages. Thus, traffic from an invalid source is encapsulated and forwarded by a first hop router (consuming its resources), and it is not dropped until it reaches the RP (consuming network and RP resources, and providing an opportunity DOS attacks). In Bi-directional PIM, these problems are even more aggravated by the fact that a common (*, G) forwarding entry is used to forward traffic in both directions (i.e., upstream towards the RP, and downstream away from the RP and towards the receivers). Although there is no protection is available to protect against malicious or unwanted sources, there is protection from unwanted receivers joining a multicast group through admission-control mechanisms that are available.\nOne approach to limit the problem is to use security Access Control Lists (ACLs) manually maintained by an administrator to limit the hosts which can send multicast traffic, such as by denying all but an explicit list of hosts. Clearly such a solution is not scalable."}
{"text": "The present invention relates to the field of local area networks in general and specifically to improvements in managing network traffic between local area network segments and the network backbone.\nLocal area networks are communication networks extending from several hundred to several thousand feet within a geographical region. These networks provide efficient and high-speed communications between host computers and other devices for the purposes of sharing resources and communicating in a distributed processing environment. For example, local area print servers connected to a local area network eliminate the need for dedicating a printer to each computer. Local area file and disk servers permit the shared use of large data bases without the cost of maintaining multiple copies of the data bases.\nOne type of local area network that has gained widespread acceptance is the Ethernet local area network. Ethernet is a carrier sense/multiple access with collision detection system (CSMA/CD) utilizing coaxial cable. This system is described in an article by P. M. Metcalfe and D. R. Boggs in Communications of the ACM, July, 1976, and is further described in IEEE specification 802. The Ethernet network transmits information between devices connected to the network in data packets or frames that contain a destination address corresponding to one or more of the devices.\nA local area network such as Ethernet may connect several hundred or even thousands of devices or nodes. As the number of nodes on the local area network increases, and the traffic on the network also increases, the network may be divided into smaller segments in order to manage the network traffic more efficiently. Specifically, one or more local network segments may be created to isolate local message traffic from the rest of the network. This is particularly useful where, for example, the local traffic consists primarily of data packets transferred between the nodes on the local network segment and a disk server or other type of host computer.\nThe local segments of a network are connected to a network backbone for transferring data packets between other nodes on the network and other local segments. Networks and network segments are typically connected by using bridges. The function of a bridge is to receive data packets from one network or network segment and to forward those data packets to the other network or network segment if the data packet contains a destination address corresponding to a node on that other network or network segment. Thus, bridges must act as perfect filters between the networks.\nA network bridge, however, is an expensive and complex device. In order for the bridge to filter all data packets transferred on the networks connected by the bridge perfectly and thus completely isolate the traffic on one network or network segment from the rest of the network and vice versa, the bridge must store large tables of addresses corresponding to the nodes on the networks. Furthermore, if the number of nodes on the overall network is large, it becomes impractical to filter all data packets from the network backbone to the local network segment perfectly."}
{"text": "Not applicable.\nThe present invention relates to the general field of telephone and face-to-face communications. In more particular, the invention relates to systems to assist telephone and face-to-face communications by those persons who are deaf, hard of hearing, or otherwise have impaired capability to receive or understand communications that are transmitted by voice.\nMost modem human communications in both social and business environments takes place through sound communications. Yet within modern society there are many persons who have attenuated hearing capability. To assist those persons in making use of our telephonic communication system built for the hearing majority, there has been developed a system of telephone communication which has been principally used by the deaf community. That system makes use of a category of device known variously as a telecommunication device for the deaf (TDD), text telephone (TT) or teletype (TTY). Current TDDs are electronic devices consisting of a key board and a display as well as a specific type of modem, to acoustically or directly couple to the telephone line. Modem TDDs permit the user to type characters into their keyboard, with the character strings then encoded and transmitted over the telephone line to be displayed on the display of a communicating or remote TDD device.\nMost TDD communication is conducted in an idiosyncratic code specific to the community of TDD users. This code, known as Baudot, evolved historically at a time when many telelecommunication devices for the deaf were based on mechanical or electromechanical devices rather than the current technology based on digital electronic components. Accordingly, the Baudot protocol was constructed for a set of constraints which are no longer relevant to present date devices. The original Baudot protocol was a unidirectional or simplex system of communication conducted at 45.5 Baud. The conventional Baudot character set was a character set consisting of 5 bit characters and the system encodes the bits of those characters in a two-tonal system based on carrier tones of 1400 and 1800 Hertz.\nThe system of TDD communications is widely used and in fact has become indispensable to the deaf community throughout the industrialized world. Deaf persons extensively communicate with their neighbors and with other deaf and hearing people remotely, using the TDD system. In addition, systems have been developed to facilitate the exchange of communication between the deaf community and hearing users who do not have access to or utilize a TDD device. In the United States, telephone companies have set up a service referred to as a xe2x80x9crelay.xe2x80x9d A relay, as the term is used herein, refers to a system of voice to TDD communication in which an operator, referred to as a xe2x80x9ccall assistant,xe2x80x9d serves as a human intermediary between a hearing user and a deaf person. Normally the call assistant wears a headset that communicates by voice with the hearing user and also has access to a TDD device which can communicate to the deaf user using a TDD appropriate protocol. In normal relay operations in the prior art, the call assistant types at a TDD keyboard the words which are voiced to her by the hearing user and then voices to the hearing user the words that the call assistant sees upon the display of his or her TDD. The call assistant serves, in essence, as an interpreting intermediary between the deaf person and the hearing person to translate from voice to digital electronic forms of communication.\nA limitation in the effectiveness of current relay protocols is the necessity for the call assistant simply to type what is said. Typical call assistants can usually type reasonably well, typically in the range of 30 to 50 words per minute. While systems exist which permit the digital encoding of verbal communications at a faster rate, such as court stenography used in the legal system, such systems require extensive special training and are impractical for the numerous call assistants required by the relay systems in operation today. This limitation on the speed of conversion from speech to digital communications hampers the effective flow of communication in a relay call between the hearing person and a deaf person, since most hearing people speak at a rate which is higher than most call assistants can type. In addition, since conventional Baudot communications is unidirectional, the flow of conversation in a relay assisted communication session can be somewhat awkward. For example, first the hearing person must voice a statement or question. Then the call assistant must type that statement or question, which is then transmitted at Baudot speeds, which are slower than normal human voice communication, to the deaf person. The deaf person waits until the entire statement or question is transmitted to him or her, after which he or she composes a response and types it in at his or her TDD. Then the communication flows backward to the call assistant who must voice to the hearing person what the deaf person has typed at his or her terminal. This process enables a degree of two-way communication between a deaf person and a hearing person, but the system tries the patience of the hearing person, since it is typically not conducted at a pace anywhere close to normal human communications.\nThe present invention is summarized in that a relay system to facilitate the translation of information and communication between deaf and hearing persons includes a call assistant who re-voices the words of the hearing person which are spoken to the call assistant. The words spoken by the call assistant are recognized by a speech recognition computer program which has been trained to the voice pattern of the call assistant, such that the words are promptly translated into a high speed digital communication protocol. That high speed digital communication message is then transmitted electronically promptly by telephone to a visual display accessible to the deaf person.\nIt is an advantage of the invention described herein that the call assistant does not have to type most, if any, of the words spoken by the hearing person in the communication session so that the overall speed of communications from the hearing person to the deaf person is dramatically increased.\nIt is an object of the present invention that the design and utilization of a relay operated in accordance with the protocols described herein permits the introduction of small hand-held personal interpreter which will enable on the spot communications between an assisted user and a hearing person wherever the assisted user might be.\nOther objects, advantages and features of the present invention will become apparent from the following specification when taken in conjunction with the accompanying drawings."}
{"text": "1. Field\nThis is related to a semiconductor device having a function, for limiting (cutting) the power supply of an internal circuit.\n2. Description of the Related Art\nInclusion of a function of partially limiting (cutting) the power supply of an internal circuit, namely, a so-called power gating function, in large-scale integrated circuits (LSIs) to decrease the power consumption is becoming the mainstream. On the other hand, LSIs have a plurality of test circuits for use in a test of the LSIs. More specifically, the LSIs have circuits, such as an input/output test circuit that performs fault diagnosis of data writing/reading circuits and an accelerated test circuit that performs an accelerated test on internal circuits of the LSIs.\nAn input/output test circuit includes a circuit that serially couples, in series like a chain, flip flops related to writing and reading of an LSI and a control circuit that receives control signals and test data from outside of the LSI and controls the state of the test (see, for example, Japanese Laid-open Patent Publication No. 2003-98223). In addition, an accelerated test circuit includes a circuit that applies an accelerated test voltage to an internal circuit and a control circuit that receives control signals from outside of an LSI and controls the state of the test (see, for example, Japanese Laid-open Patent Publication No. 10-21699).\nA test of an LSI having a control circuit that performs a power gating operation and a control circuit that controls the state of the test has to be carried out while allowing the control circuits to operate in cooperation with each other. However, sending a control signal for permitting the cooperation between the control circuits to an LSI from a testing equipment residing outside of the LSI requires an expensive testing equipment having an advanced function, which thus increases the cost for the test.\nNevertheless, when the expensive testing equipment that generates a control signal for permitting the cooperation between control circuits does not employed and the control circuit that performs the power gating operation and the control circuit that controls the state of the test do not operate in cooperation with each other, the following disadvantages may occur. At the time of an accelerated test, application of a sufficient level of stress, which is carried out by applying a predetermined voltage level to an internal circuit, is not achieved because the power gating function is not cancelled. Additionally, at the time of an input/output test, test data is not transferred since an unexpected power gating operation is activated. As a result, the accelerated test or the input/output test may be undesirably prevented."}
{"text": "Known groomer vehicles are normally powered crawler vehicles comprising a frame extending along a first axis; and a snow grooming implement which is normally a rotary snow tiller drawn by the groomer vehicle and comprising a second frame extending predominantly crosswise to and symmetrically with respect to a second axis which, when traveling along a straight path, is aligned with the first axis, and, when traveling along a curved path, forms an angle with the first axis.\nSuch known hitch devices comprise an arm or so-called boom hinged to the first frame about a third axis perpendicular to the first axis, and about a fourth axis perpendicular to the first and third axis; a first actuator is connected to the first frame and the arm to adjust the position of the arm about the third axis with respect to the first frame; and at least one second actuator is connected to the first frame and the arm to adjust the position of the arm about the fourth axis.\nThe first actuator is often a hydraulic cylinder for lifting and lowering the arm to selectively lift and lower the rotary snow tiller with respect to the snow covering. The second actuator sets the angle between the first and second axis: when traveling along a straight path, the second actuator keeps the arm aligned with the first and second axis; and, when traveling along a curved path, the second actuator sets the arm to form a given angle with the first axis about the fourth axis.\nKnown hitch devices have proved highly effective, both in lifting and lowering the rotary snow tiller, and as regards maneuvering around bends and stabilizing the position of the tiller.\nIn some ski slope grooming conditions, however, it is necessary for the rotary snow tiller to go over an already groomed strip of the snow covering without the groomer vehicle traveling over it.\nKnown hitch devices allow this to be done, but force the rotary snow tiller to operate in a skew position with respect to the traveling direction, which increases energy consumption and reduces grooming width.\nMoreover, when traveling along a curved path, known hitch devices are not always able to position the rotary snow tiller to cover the prints left by the groomer vehicle in the snow covering."}
{"text": "This invention relates generally to the field of child products, and in particular to bibs. More specifically, the invention relates to the use of bibs that are attached to neck pillows.\nBibs have long been used to protect a child from becoming soiled. Bibs are most often used when a child is eating, although they may be used for other applications as well. Bibs are typically constructed of a piece of material that is coupled about the child's neck using a pair of strings or other fasteners, such as snaps of a hook and loop fastener material.\nOne particularly useful bib is described in U.S. Pat. No. 6,523,200. This invention is related to other types of bibs that may be used with neck pillows."}
{"text": "Several Doppler phantoms have been described in the prior art. These fall into two main categories.\nIn the first group are phantoms which include a moving target such as a ball bearing or suspended string. The moving target is used to map the location and size of the sensitive Doppler sample volume of the ultrasound device. In this category, for example, Walker et al, \"Evaluating Doppler Devices Using a String Target\", Journal of Clinical Ultrasound, 10:25, 1982, describe a method in which a pair of moving strings at two different depths is used in a test target to measure Doppler sample volume size, sample location, and amplitude sensitivity to Doppler shift. An alternative method, described by Hoeks et al, \"Methods to Evaluate the Sample Volume of Pulsed Doppler Systems\", Ultrasound in Medicine and Biology, 10:427, 1984, is to couple a small target such as a sphere to a vibrating loudspeaker and then map the Doppler sample volume using this moving target.\nIn the second category of Doppler phantoms, the objective is to simulate blood flow, using a blood-mimicking liquid, through a simulated blood vessel in a tissue-mimicking material. Phantoms of this kind typically consist more or less of a solid block of tissue-mimicking material, such as gelatin or polymer, which contain tubes of varying diameters and bifurcations. Such phantoms seek to simulate blood flow through abdominal and peripheral vessels. Examples of such devices, of varying degrees of sophistication, are taught by Newhouse et al, \"A Proposed Standard Target for Ultrasound Doppler Gain Calibration\", Ultrasound in Medicine and Biology, 8, 313--1982; McDicken, \"A Versatile Test-Object for the Calibration of Ultrasonic Doppler Flow Instruments\", Ultrasound in Medicine and Biology, 12:245, 1986; and Boote et al, \"Performance Tests of Doppler Ultrasound Equipment with a Tissue and Blood Mimicking Phantom\", Journal of Ultrasound Medicine, 7, 137-147, 1988. phantoms allow independent measurement of \"blood\" flow and enable calibration of the fluid velocity estimations of Doppler ultrasound. Commercial versions of these test objects, \"Tissue Mimicking Ultrasound Phantom \" Model 409, Radiation Measurements, Inc., Middleton, Wis. 53362 (U.S. Pat. No. 4,277,367); \"Ultrasound Doppler Phantom\", ATS Laboratories, Inc., Bridgeport, Conn. 06608; and \"Ultrasound Doppler Phantom\", Interspec, Inc., Lewiston, Pa. 17044, are also available.\nThe most extensive application of medical Doppler ultrasound, however, is in cardiac diagnosis. Doppler ultrasound is used to measure blood flow through cardiac valves and thus to estimate pressure drops across individual heart valves for detection of valve dysfunctions. In addition, cardiac Doppler ultrasound is used to detect anatomical anomalies such as ventricular septal defects. The field of cardiac ultrasound imaging and Doppler ultrasound has a known and long-standing need for an anthropomorphic cardiac Doppler phantom to assess the performance of Doppler ultrasound under more realistic clinical conditions. Limited attempts to develop a Doppler phantom for cardiac applications include a modified cardiac pulse duplicator, \"Cardiac Pulse Duplicator\", Model MP1, Dynatek Laboratories, Annandale, N.J. 08801 as described by Cary et al in a private communication. However, this device is fabricated from rigid plastic, which generates strong acoustic reverberations unsuitable for diagnostic ultrasound examinations. In addition, such a modified pulse duplicator shows no similarity to the contracting muscular heart of human anatomy. Other examples include a modified cardiac pulse duplicator, \"Valve Visualization Pulse Duplicator\", Model MV/T1 Dynatek Laboratories, Annandale, N.J. 08801, as described by Gels et al, \"In Vitro Ultrasound Flow Imaging Through Prosthetic Heart Valves\", Medical Instrumentation, 21(2), 66-74, 1987, that includes a soft rubber vessel to simulate the aorta downstream of the aortic valve. This system, using microbubbles in tap water as an ultrasound contrast agent, enables flow visualization with diagnostic ultrasound imaging equipment to evaluate, qualitatively, the performance of prosthetic heart valves in vitro but does not permit simulation of the contracting heart of human anatomy.\nIn addition, Reul, \"Cardiovascular Simulation Models\", Life Support Systems, 2, 77-98, 1984, has developed hydraulically driven cardiovascular simulation models for the evaluation of prosthetic heart valves, using pressure transducers, flow transducers, and optical video filming of suspended particles in an aqueous-glycerol solution, but the cardiac models themselves are housed in rigid polymethylmethacrylate boxes. There are no viewing ports in the containers for ultrasound imaging or ultrasound Doppler studies. The polymethylmethacrylate containers produce high attenuation and reverberations that are unsuitable for ultrasound imaging or ultrasound Doppler studies. The hydraulic drive unit is mounted below the left ventricle in this device and precludes an apical ultrasound view, one of the most common views in clinical cardiac diagnosis. The hydraulic medium surrounding the left ventricle also is not specified. However, the attenuation scattering and velocity values of this medium are critical to the development of an ultrasound phantom as will be described below. Ordinary or distilled water, for example, is not suitable.\nThe field of cardiac assist devices, furthermore, includes left ventricular assist devices (LVAD) that typically have pneumatic pumps composed of a rigid plastic shell surrounding the flexible polymer left ventricle, Van Citters et al, \"Artificial Heart and Assist Devices: Directions, Needs, Costs, Societal and Ethical Issues, \" Artificial Organs, 9(4), 375-415, 1985. The presence of air around the blood sac and the hard plastic shell make this technology unsuitable for a cardiac ultrasound phantom. One hydraulically driven LVAD is known which also consists of a rigid case surrounding the left ventricle, Altieri et al, \"Implantable Ventricular Assist Systems\", Artificial Organs, 11(3), 237-246, 1987. Again, however, this shell precludes use of the LVAD as an ultrasound phantom.\nA need clearly exists in the art of Doppler and ultrasound imaging and the like for medical purposes for an anthropomorphic ultrasound phantom that realistically enables ultrasound analyses of simulated defects in human hearts and provides reverberation free outputs suitable for comparison of different devices and simulation studies."}
{"text": "The present invention relates to adhesively bonded engine component assemblies and more particularly to adhesively bonded intake manifold assemblies. The invention also relates to an intake manifold assembly employing a cure-on-demand adhesive.\nHistorically, fabrication techniques have involved the fabrication of separate metal components and joining the components using fasteners, with an intermediate disposed gasket. This tends to require many parts and therefore is labor intensive.\nPlastic engine intake manifold assemblies (such as for automotive vehicle engines) have been fabricated by near net shape manufacturing processes, such as lost core processes. While these tend to result in desirable near net shape products, the shapes of which may be complex, they typically are expensive.\nIn another popular process, two or more plastic intake manifold components are made, each with a planar mating surface. The components are joined at each of their respective mating surfaces by vibration welding. This process also tends to be time, energy and labor intensive. Further, there is a potential that any warpage of respective components or other surface irregularities will result in undesirable gaps at the mating surface upon welding. Thus, it is imperative that opposing mating surfaces are planar, which constrains design flexibility.\nAccordingly, there is a need to provide a method of manufacture and resulting assembly, pursuant to which multiple components of an engine intake manifold assembly are joined together with an adhesive bond. There is a further need for respective components to be attached to each other in the optional absence of fasteners, a gasket or both. There is still a further need for a high integrity joint be achieved opposing mating surfaces that are non-planar, thereby enabling intricate plastic component shapes (with or without additional integrated components) to be molded by conventional molding processes.\nThe present invention meets the above needs and others by providing an engine intake manifold assembly including a first component having a first mating surface, and a second component having a second mating surface. A layer of adhesive is provided between the first and second components in contact with the respective mating surfaces for joining the first and second components to define an engine intake manifold assembly."}
{"text": "1. Technical Field\nThe present disclosure relates to electronic devices and, particularly, to an electronic device with a dummy connector.\n2. Description of Related Art\nReferring to FIGS. 4 and 5, a computer system 300 includes a shell 310 and a printed circuit board 320 accommodated therein. The printed circuit board 320 includes a plurality of connectors 322, each fixed on the printed circuit board 320 by inserting pins of the connector 322 into corresponding fixing holes 324 defined in the printed circuit board 320, and welding the pins to the printed circuit board 320 by solder. The shell 310 defines a plurality of through holes 312. The plurality of connectors 322 are exposed respectively through the plurality of through holes 312 to connect to corresponding system elements, enabling communication therewith.\nWith computer systems available in many different configurations, locations of connectors can change, and in certain cases even removal of some is necessary. If one of the plurality of connectors 322 is not used, a corresponding through hole 312 in the shell 310 should be eliminated as well to avoid electromagnetic leakage therethrough. Accordingly, production process of the shell 310 needs to be adjusted. However, such an adjustment is complicated and costly."}
{"text": "Control is difficult in particle accelerators when bunches with different characteristics have to travel through the same beam pipe and see the same optics. Controlling bunches of different types with the same magnets is a challenge. Optimizing the lattice for one type of bunch typically worsens control for the other types of bunches. Unfortunately, no knobs are available that affect only one type of bunch while leaving the others untouched. As an example, in a recirculator energy recovery linacs (ERL) having a single linac, the accelerated and decelerated beams have the same energy and thus cannot be separated by a spreader/recombiner.\nAs an example, JLAMP (JLab AMPlifier), a 4th generation light source covering the range 10 eV-100 eV in the fundamental mode with harmonics to 1 keV, includes common recirculation transport of accelerated and recovered beams. The common recirculation transport imposes significant constraints on both performance and operability, as multiple beams must be managed using a common set of accelerator components. Common transport is therefore a significant problem in multi-pass ERL designs having a single linac.\nAccordingly, what is needed is a method for achieving independent control over the bunch types in common transport lines. As an example, independent control over the bunch types could be used to increase the power of an extreme ultra violet (EUV) free electron laser (FEL), which would be of interest in lithography systems, or in ERL-based light sources such as JLAMP, eRHIC, LHeC, CERN, and KEK."}
{"text": "(a) Technical Field\nThe present invention relates to an induction heating device for a fuel cell system. More particularly, it relates to an induction heating device for a fuel cell system, which can to rapidly heat coolant during cold start-up, control the power consumption depending on the voltage of a fuel cell stack, and ensure the insulation resistance by separating a heating unit, which is in contact with the coolant, from the outside.\n(b) Background Art\nA conventional fuel cell system comprises a fuel cell stack for generating electricity and heat in the form of reaction products via an electrochemical reaction. In a fuel cell system, the chemical energy of a fuel is converted into electrically energy. Typically, fuel cell systems include, for example, a hydrogen supply system for supplying hydrogen as fuel to the fuel cell stack, an oxygen (air) supply system for supplying oxygen-containing air as an oxidant required for the electrochemical reaction in the fuel cell stack, and a thermal management system (TMS) for removing reaction heat from the fuel cell stack to the outside of the fuel cell system, controlling operation temperature of the fuel cell stack, and performing water management function.\nA fuel cell stack is consists of a plurality of unit cells. Each of these unit cells have an anode, a cathode and an electrolyte (electrolyte membrane). Hydrogen is supplied to the anode (“fuel electrode”) and oxygen-containing air is supplied to the cathode (“air electrode” or “oxygen electrode”).\nThe hydrogen supplied to the anode is dissociated into hydrogen ions (protons, H+) and electrons (e−) by a catalyst disposed in an electrode/catalyst layer. The hydrogen ions and electrons are transmitted to the cathode through the electrolyte membrane and a separator, to respectively. At the cathode, the hydrogen ions and electrons transmitted from the anode react with the oxygen-containing air supplied to the cathode to produce water. Electrical energy generated by the flow of the electrons is supplied to a load (e.g., a motor for driving a fuel cell vehicle) requiring the electrical energy through a current collector of an end plate.\nThe role of the thermal management system for the fuel cell system is particularly important in order to prevent the temperature of the fuel cell stack from increasing during operation of the fuel cell system and improve cold startability.\nThe conventional configuration of a thermal management system will be described with reference to FIG. 4 below. In FIG. 4, a water pump 16 is provided in a coolant circulation line 14, which is connected between a coolant outlet of a fuel cell stack 10 and an inlet of a radiator 18. The conventional thermal management system also has a coolant heating device 20 and a reservoir which are sequentially connected in the coolant circulation line 14.\nHere, the coolant circulating through the coolant circulation line 14 functions as a cooling medium that cools the fuel cell stack 10 during operation and functions as a heating medium that is rapidly heated by the coolant heating device 20 and supplied to the fuel cell stack 10 to rapidly thaw the fuel cell stack 10 during cold start-up.\nIn conventional designs, the coolant heating device 20 includes a heater for heating the coolant. The heater is used to increase the temperature of the coolant, achieve more rapid temperature stabilization of the coolant during start-up of the vehicle, and reduce the to remaining voltage of the fuel cell stack during shutdown of the vehicle.\nIn FIG. 4, the conventional coolant heating device 20 has a structure, in which cartridge resistance heaters 24 are connected to a power unit provided in a housing 22, through which the coolant passes. However, this structure has the following disadvantages. First, there should be a predetermined distance between each of the heaters 24, which are high density heating elements. Thus, the volume of the housing 22 is significantly increased in conventional designs. Therefore, it is difficult to ensure a space for mounting the heating device which has such a large volume in a compact vehicle such as a fuel cell vehicle. Moreover, the complicated structure increases the manufacturing cost and reduces the assemblability and maintainability of this system. Furthermore, the cartridge resistance heaters are in direct contact with the coolant, and thus it is necessary to ensure the cartridge resistance heaters are insulated accordingly.\nIn addition, the cartridge resistance heaters penetrating the housing are typically sealed with O-rings, and thus when the temperature of the heaters is increased, the coolant may leak due to deformation of the O-rings, thus destroying the insulation.\nFIG. 5 is a schematic diagram showing a conventional cold start-up system for a fuel cell. Referring to FIG. 5, a plurality of resistors R1, R2 and R3 and a plurality of switches L1, L2 and L3 are connected to a fuel cell stack 10 such that the resistance values are controlled step by step using the resistors R1, R2 and R3. As a result, it is possible to prevent the voltage of the fuel cell stack from dropping below a lower limit, thus enabling the cold start-up.\nIn detail, the plurality of resistors R1, R2 and R3 having different resistance values are connected in parallel to both ends of the fuel cell stack 10, and the plurality of switches L1, L2 and L3 for switching the circuit are connected to the resistors R1, R2 and R3, respectively, to change the resistance values applied to the both ends of the fuel cell stack 10 by on/off control of the switches L1, L2 and L3 such that the voltage of the fuel cell stack 10 is controlled between a predetermined lower limit for achieving cold start-up and an upper limit during cold start-up.\nHowever, to control the power of the fuel cell stack during cold start-up, this system requires a heater at each resistor and a relay for controlling the heater, which imposes various limitations to the overall system.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "The present invention relates to parachutes and more particularly low altitude parachutes for use in both high velocity deployments as well as low velocity deployments.\nThe recovery of personnel resulting from low altitude ejections from military aircraft has become an increasingly difficult task. For many military applications the operational velocities are becoming greater causing the recoverable payloads such as ejection seats and ejectable capsules to become heavier due to the addition of protective equipment. The heavier payloads typically require larger parachutes so that they may be successfully recovered. The larger parachutes, however, take longer to inflate which may result in premature ground impact if deployed at low speeds and low attitudes.\nParachute opening speed is proportional to its deployment speed. This generates two major concerns which the recovery system must address. The two major concerns occur at the velocity extremes of the operational envelope. The first concern, as discussed above, is premature ground impact at low speeds due primarily to slow inflation. The second concern is the excessive opening forces that are present at high speeds.\nAt low speeds, the airflow surrounding the parachute is often too slow to fully inflate a large parachute in a relatively short time. This is due primarily to the inherently poor opening characteristic of standard parachutes which wastes or does not effectively utilize the available inflation air. There is often very limited time in which to inflate the parachute before ground impact. Thus in the low speed scenario, a major design goal is to speed up the inflation process of the parachute by capturing much of the available inflation air.\nThe opposite is true in the high speed regime, where the airflow velocity is sufficient to quickly inflate the parachute. The concern, however, is that the drag force caused by the parachute is proportional to the square of the deployment velocity multiplied by the effective drag area of the inflated parachute. Thus at high speeds when the parachute inflates very quickly, the drag forces can become excessive. If not controlled, these drag forces can exceed human tolerance limits or otherwise cause failures in the parachute, which is often equally fatal to the crew members. Therefore in the high speed regime, the design goal is to slow down or control the inflation process of the parachute in order to avoid excessive drag forces. Alternatively, one can attempt to control the drag forces by controlling the effective drag area so as not to exceed tolerable limits.\nThe related technology for speeding up the inflation of parachutes include a variety of techniques. The conventional methods for speeding up the inflation of a parachute include; reducing the permeability of the fabric or the porosity of the canopy or both. Also, one can close the vent of the parachute which prevents the trapped air from escaping. Yet another conventional method for speeding up the inflation of a parachute is to sew inflation pockets on the canopy around the outside of the skirt. These inflation pockets aid the spreading of the canopy by generating greater eddies of air around the parachute and contributing a small spreading force. These conventional methods, however, do not increase the opening speed of the parachute enough to warrant their use on low altitude military ejection applications.\nOne method currently used on some ejection seat parachutes to speed up inflation at low speed is a pull down vent line. This technique is also commonly known as pulling down the apex of the parachute. In this method a line attached to the vent area of the parachute is pulled down prior to deployment. When the parachute is deployed at low speed, the pulled-down vent forces the air entering the parachute to open the skirt quickly, allowing for a faster inflation. When the parachute is deployed at high speed, the pull-down vent line simply breaks allowing for normal inflation. The increase in opening speed of the parachute at low velocities is accomplished because the opening shape of the parachute is altered. However, by altering the shape of the parachute, the distribution of forces on the canopy are also altered. Altering the opening loads presents additional problems to those involved in parachute design.\nAnother way of speeding up inflation sufficient for use on ejection seat parachutes is with the use of a spreading gun. This gun explosively throws outward metal slugs that are attached around the perimeter of the parachute skirt. This spreading force acts to open the skirt very quickly. At high velocities this fast opening of the full canopy is undesirable because it can cause excessively high opening forces.\nA third method used in this field for speeding up inflation is to use a cluster of smaller parachutes instead of one large parachute. While it is true that smaller parachutes will inflate faster, the increase in weight and complexity of such a system is viewed as a severe drawback.\nOn the other hand, the related technology for delaying the inflation of a parachute so as to avoid excessive opening forces is staged deployment systems. One such staged deployment system is a reefed parachute. This system employs a reefing line that prevents the parachute skirt from fully opening until a timed pyrotechnic line cutter severs the reefing line allowing the skirt to fully open. This type of staged deployment system distributes the opening forces in two or more smaller peak forces instead of one large peak force. Reefed parachutes and other staged deployment systems are not typically desireable during low velocity deployments. This is because any delay in inflation of the parachute can result in premature ground impact.\nThe present invention is a single parachute system which can recover a payload throughout virtually the entire performance envelope of military aircraft. That is to say, the regulated area parachute disclosed herein is capable of safely and successfully opening in both low speed situations as well as the high speed situations. The present invention is a significant advancement in recovery system technology which addresses the age old problem of minimizing the parachute inflation time while also minimizing the opening forces. This new parachute is especially effective in solving the problem of low altitude recovery of manned ejection seats and ejectable capsules."}
{"text": "1. Field of the Invention\nThe present invention relates to a PC Card input/output device for use with PC Cards and a PC Card connector, which are used for writing and reading data to and from a PC Card containing a number of storage elements, and, more particularly, to improvements in the convenience of PC Cards.\n2. Description of the Prior Art\nReferring next to FIG. 7, there is illustrated an external view of the outward appearance of a PC Card. In the figure, reference numeral 1 denotes a PC Card, 6 denotes a female connector formed on the front side of the PC Card 1 which is to be in the lead when inserted into a PC Card socket, and comprised of a plurality of female connecting terminals, 7 denotes a notch formed on the right (or upper) side and one lateral side of the PC Card 1 and extending along the length of the PC Card 1 from the front side of the PC Card 1 and, and 8 denotes a guide slit formed on the other lateral side of the PC Card 1 and extending along the length of the PC Card 1 from the front side of PC Card 1.\nReferring next to FIG. 8, there is illustrated a block diagram showing the structure of a prior art input/output device for use with PC Cards and a peripheral circuit. In the figure, reference numeral 2 denotes a PC Card connecter into which the PC Card 1 gets inserted, having a male connector which is electrically connected with the female connector 6 of the PC Card 1, and 3 denotes an internal circuit which can be electrically connected with the PC Card 1 by way of the PC Card connector 2. Referring next to FIG. 9, there is illustrated a perspective view showing the structure of the prior art PC Card connector 2 shown in FIG. 8. In FIG. 9, reference numeral 9 denotes a socket portion into which the PC Card 1 gets inserted so that the female connector 6 of the PC Card 1 is electrically connected to the PC Card connector 2, 10 denotes a plurality of male connecting terminals projected from an inner wall of the PC Card connector 2 at the back of the socket 9, 11 denotes a guide rail which can be engaged with the guide slit 8 of the PC Card 1 so as to guide the guide slit 8 when the PC Card 1 gets inserted into the socket portion 9, and 19 denotes an insertion blocking member which can be engaged with the notch 7 of the PC Card 1 when the PC Card 1 gets inserted into the socket so as to prevent the PC Card 1 from being pressed against the inner wall of the PC Card connector 2 after the plurality of female connecting terminals of the PC Card 1 are surely engaged with the plurality of male connecting terminals 10 of the PC Card connector 2, respectively.\nIn operation, the PC Card 1 gets inserted into the socket portion 9 with the upper side of the PC Card 1 pointing upwards while the guide slit 8 of the PC Card 1 is engaged with the guide rail 11. Once the PC Card 1 is inserted into the socket portion 9 as far forward as possible until the notch 7 is engaged with the insertion blocking member 19, the PC Card 1 is stopped and the plurality of female connecting terminals of the connector 6 of the PC Card 1 are brought into contact with and therefore electrically connected to the plurality of male connecting terminals 10 of the PC Card connector 2, respectively.\nAs a result, the internal circuit 3 is connected to the plurality of female connecting terminals of the PC Card 1 by way of the PC Card connector 2, so that the PC Card 1 can furnish and receive data to and from the internal circuit 3 in response to reading and writing requests from the internal circuit 3.\nIt is clear from FIG. 9 that the PC Card 1 cannot get inserted into the socket portion 9 with the reverse or lower side of the PC Card 1 pointing upwards because the front edge of the PC Card 1 is brought into contact with the insertion blocking member 19 and therefore the insertion of the PC Card 1 is stopped on the way to the plurality of male connecting terminals 10. The plurality of female connecting terminals of the PC Card 1 thus cannot reach the plurality of male connecting terminals 10 of the PC Card connector 2.\nAs previously explained, when utilizing a PC Card 1 and a PC Card connector 2 which comply with the PC Card Standard, the PC Card 1 can get inserted into the socket portion 9 of the PC Card connector 2 only when making the upper side of the PC Card 1 point upwards. As a result, the PC Card 1 can be connected to the internal circuit while a certain connection between the PC Card 1 and the internal circuit is established.\nWhile the combination of the above-mentioned prior art PC Card and the PC Card connector can prevent the user from inserting the PC Card while mistaking the lower side of the PC Card for the upper side of the PC Card, and therefore prevent improper electrical connections between the PC Card and the PC Card connector, there is a problem in that every time the user does an improper insertion the user has to insert the PC Card again by turning the PC Card upside down, resulting in reduction in the convenience of the PC Card. Especially, when the PC Card connecter is disposed such that its insertion opening is running longitudinally, it is difficult for the user to determine in which direction the upper and lower sides of the PC Card should be oriented, resulting in substantial reduction in the convenience of the PC Card.\nJapanese Patent Application Laying Open (KOKAI) No. 6-155974 discloses an improved prior art PC Card to solve the above problem. Referring next to FIG. 10, there is illustrated a perspective view showing the outward appearance of the PC Card as disclosed in the reference. In FIG. 10, reference numeral 20 denotes a switch disposed on one lateral side of the improved PC Card 100 which is opposite to the other lateral side on which a notch 7 is formed. The switch 20 can slide on the lateral side of the PC Card 100 opposite in direction to the travel of the PC Card 100. The switch 20 extends such that its free end reaches the front side of the PC Card 100 which is to be in the lead when the PC Card 100 gets inserted into a PC Card socket. Reference numeral 21 denotes a connection switching unit connected between the female connector 6 and a memory disposed inside the PC Card 100, responsive to a switching signal from the switch 20, for changing the connection between the plurality of female connecting terminals of the PC Card 100 and the pins of the built-in memory. To be more specific, the connection switching unit 21 performs a switching in such a manner that the connecting target with which each of the pins of the built-in memory is to be connected is switched from one female connecting terminal of the PC Card 100 to another female connecting terminal symmetric to the former female connecting terminal with respect to the center C of the front side of the PC Card 100.\nWhen the user inserts the PC Card 100 into the PC Card connector with the upper side of the PC Card pointing upwards, the PC Card 100 is properly mounted into the PC Card connector, like the above-mentioned prior art PC Card as shown in FIG. 9. On the other hand, when the user inserts the PC Card 100 into the PC Card connector with the reverse side of the PC Card pointing upwards, the free end of the switch 20 is brought into contact with the insertion blocking member 19. When the user further pushes the PC Card 100 into the back of the socket, the female connector 6 of the PC Card 100 is electrically connected to the PC Card connector 2. In the PC Card 100, the switch 20 then generates a switching signal in response to the travel of the switch 20, with the result that the connection switching unit 20 changes the connection between the plurality of female connecting terminals and the pins of the built-in memory in such a manner that the connecting target with which each of the pins of the built-in memory is to be connected is switched from one female connecting terminal to another female connecting terminal symmetric to the former female connecting terminal with respect to the center C of the front side of the PC Card 100.\nAccordingly, even though such the prior art PC Card 100 is inserted improperly into a PC Card socket with the reverse side of the PC Card 100 pointing upwards, the improved PC Card 100 can be electrically connected to the PC Card connector while the connection between the built-in memory and the internal circuit 3 by way of the female connector 6 and PC Card connector 2 is established in the same manner as the case where the PC card 100 is properly inserted into the PC Card socket with the right side of the PC Card 100 pointing upwards because the electrical coupling between the plurality of female connecting terminals of the female connector 6 of the PC card 100 and the plurality of male connecting terminals 10 of the PC Card connector 2 has been changed and the connection between the plurality of female connection terminals of the PC Card 100 and the pins of the built-in memory has been changed similarly. Thus the PC Card 100 can receive and furnish data from and to the internal circuit 3 reliably in response to reading and writing requests from the internal circuit 3 regardless of how the PC Card 100 inserted into the PC Card socket is oriented.\nWhile the improved PC Card enabling the improper insertion offers an advantage of allowing an internal circuit to read and write data from and to the PC Card inserted into a PC Card socket regardless of how the PC Card inserted is oriented, it cannot improve the operability of prior art PC Cards which have become widely available and which complies with the PC Card Standard, instead, there is a fear that the introduction of the improved PC Card on the market causes a state of disorder wherein the improved PC Cards enabling the improper insertion coexist with existing prior art PC Cards and therefore users are required to use both the improved PC Cards and prior art PC Cards each in its proper way, resulting in reduction in the general convenience of PC Cards."}
{"text": "Amplitude-shift-keyed modulation is a commonly known technique used to modulate rf and optical carriers for the transmission of data. A typical transmission involves sending bursts of carrier signals, the presence of a burst identifying a binary one or zero and the absence of a burst or carrier being representative of the opposite binary digit. Examples of patents describing various ASK modulations are U.S. Pat. Nos. 4,947,407, 4,860,320, and 4,829,560.\nTechniques for improving detection of data signals detected with optical receivers are known in the art. In one such technique as shown and described in U.S. Pat. No. 4,431,916 delayed versions of input signals are compared with undelayed versions to reproduce data signals. Threshold values are created from these signals. U.S. Pat. No. 4,507,795 describes an apparatus for locating leading and trailing edges of pulses derived from an RF input.\nOne problem associated with ASK receivers involves the large amplitude swings of incident RF or optical carrier signals. Such large swings tend to introduce undesirable variations in the output data pulse width. Such variations often are due to the receiver's bandpass filtering which causes sloping, rising and falling edges at the output of the logarithmic amplifier that is commonly used in receivers."}
{"text": "1. Field\nEmbodiments of the invention relate to prosthetic products, and more particularly to systems and methods for decreasing build height of prosthetic products.\n2. Description of the Related Art\nMillions of individuals worldwide rely on prosthetic and/or orthotic devices to compensate for disabilities, such as amputation or debilitation, and to assist in the rehabilitation of injured limbs. Orthotic devices include external apparatuses used to support, align, prevent, protect, correct deformities of, or improve the function of movable parts of the body. Prosthetic devices include apparatuses used as artificial substitutes for a missing body part, such as an arm or leg."}
{"text": "(i) Technical Field\nThe present invention relates to image pickup devices and focal plane shutters.\n(ii) Related Art\nJapanese Unexamined Patent Application Publication Nos. 2010-181522 and 2008-256844 disclose a focal plane shutter driving a shutter by use of an electromagnet.\nIn a case of using the actuator as a drive source for driving the shutter, a change in temperature of a coil of the actuator also changes the resistance value of the coil. This might cause an exposure period to vary."}
{"text": "The present invention relates to a screen-type printing stencil having therein holes arranged in a pattern.\nScreen printing, whether of the type employing flat-bed printing stencils or of the type employing cylindrical screen printing stencils in rotary screen printing machines, requires that the boundary areas between the different printed colors be smeared as little as possible.\nThe object of the present invention is accordingly to provide a screen-type printing stencil having therein openings patterned to reduce smearing.\nThis object is achieved according to the present invention by providing that the pattern of the openings or holes is formed by repetition of a primary pattern or pattern period formed by at least three holes none of which can be brought into coincidence or superimposition by rectilinear shifting toward each other.\nIn this connection it is of advantage if the pattern period includes holes of at least two different geometrical shapes and/or different sizes.\nA particularly advantageous embodiment of a screen-type printing stencil according to the present invention is obtained when the pattern period has at least one hole which is equilateral and at least two elongated holes or six pentagonal holes.\nWhen at least two elongated holes are provided, it is advantageous if the equilateral hole is hexagonal and if there are present three elongated, congruent, rhomboidal holes, the longitudinal axes of which extend at an angle of 60.degree. with each other. The sides of each rhomboidal hole extend in directions parallel with respective sides of the hexagonal hole.\nWhen six pentagonal holes are provided, it is of advantage that the equilateral hole be a regular hexagon and that the six pentagonal holes be shifted with respect to each other by 60.degree. symmetrically around the hexagonal hole. Each side of the hexagonal hole faces a parallel side of equal length of one of the pentagonal holes. The mutually facing sides of adjacent of the pentagonal holes are always parallel with each other, and the outer sides adjacent thereto are always aligned with each other and jointly form a regular hexagon.\nThe above features of the screen-type printing stencil of the invention avoid the occurrence of intervals, which would otherwise exert a disadvantageous effect on printing, especially the printing of fine line patterns.\nFurther, a particular advantage of the invention is that the screen-type printing stencils of the invention produce sharper pattern outlines, since the inner tensions of the photolacquer exert their effect in an advantageous manner, i.e. the photolacquer does not pull back in the area of holes, which would lead to irregular edges."}
{"text": "Field\nAspects of the present disclosure relate generally to wireless communication systems, and more particularly, to physical uplink control management in long term evolution (LTE)/LTE-Advanced (LTE-A) communication systems with unlicensed spectrum.\nBackground\nWireless communication networks are widely deployed to provide various communication services such as voice, video, packet data, messaging, broadcast, and the like. These wireless networks may be multiple-access networks capable of supporting multiple users by sharing the available network resources. Such networks, which are usually multiple access networks, support communications for multiple users by sharing the available network resources. One example of such a network is the Universal Terrestrial Radio Access Network (UTRAN). The UTRAN is the radio access network (RAN) defined as a part of the Universal Mobile Telecommunications System (UMTS), a third generation (3G) mobile phone technology supported by the 3rd Generation Partnership Project (3GPP). Examples of multiple-access network formats include Code Division Multiple Access (CDMA) networks, Time Division Multiple Access (TDMA) networks, Frequency Division Multiple Access (FDMA) networks, Orthogonal FDMA (OFDMA) networks, and Single-Carrier FDMA (SC-FDMA) networks.\nA wireless communication network may include a number of base stations or node Bs that can support communication for a number of user equipments (UEs). A UE may communicate with a base station via downlink and uplink. The downlink (or forward link) refers to the communication link from the base station to the UE, and the uplink (or reverse link) refers to the communication link from the UE to the base station.\nA base station may transmit data and control information on the downlink to a UE and/or may receive data and control information on the uplink from the UE. On the downlink, a transmission from the base station may encounter interference due to transmissions from neighbor base stations or from other wireless radio frequency (RF) transmitters. On the uplink, a transmission from the UE may encounter interference from uplink transmissions of other UEs communicating with the neighbor base stations or from other wireless RF transmitters. This interference may degrade performance on both the downlink and uplink.\nAs the demand for mobile broadband access continues to increase, the possibilities of interference and congested networks grows with more UEs accessing the long-range wireless communication networks and more short-range wireless systems being deployed in communities. Research and development continue to advance the UMTS technologies not only to meet the growing demand for mobile broadband access, but to advance and enhance the user experience with mobile communications."}
{"text": "The present invention relates to a metallization technique of a semiconductor device, more specifically to a semiconductor device having a wiring or plug of a suitable structure for high integration and a method for fabrication of the semiconductor device.\nAs LSIs increase their scales, further micronization of the devices is sought. To realize semiconductor integrated circuits having gates, wirings and contact holes of more micronized dimensions, photolithography resolutions have been improved conventionally by making exposure wavelengths shorter.\nWhile the minimum resolution dimensions are thus made smaller, device structures which permit alignment margins between lithography steps to be smaller have been studied, and an attempt has been made to reduce device dimensions without reducing dimensions of patterns to be formed therein. Such device structures are exemplified by self-aligned contact (herein after called SAC), borderless contact (herein after called BLC), and so forth.\nThe conventional SAC will be explained in comparison with a structure which does not use SAC.\nIn the structure as shown in FIG. 30A in which two gate electrodes 40 are formed on a semiconductor substrate 10, and an inter-layer insulation film 20 is also formed. In opening a contact hole 22 between the two gate electrodes 40 down to the semiconductor substrate 10, it is necessary to position the gate electrodes 40, considering beforehand the alignment precision for opening the contact hole 22.\nThat is, an interval a between the contact hole 22 and the gate electrodes 40 must be secured to be at least above the alignment precision so that when a conducting film is buried into the contact hole 22, the conducting film and the gate electrodes 40 are not short-circuited (FIG. 30B). Thus, the interval between the gate electrodes 40 is dependent on the contact hole 22, which makes further micronization impossible.\nIn SAC as shown in FIG. 30C, gate electrodes 40 are covered with an inter-layer insulating film 20 and an etching selective insulating film 38. This keeps the insulating film 38 from being etched when the inter-layer insulating film 20 is etched, and when a conducting film is buried in a contact hole 22, no short-circuit takes place between the conducting film and the gate electrodes.\nAccordingly, when a misalignment takes place in a lithography step for forming the contact hole 22, the opening of the semiconductor substrate 10 is determined only by the gate electrodes 40 and the insulating film 38, which enables the gate electrodes 40 and the contact hole 22 to be arranged without considering alignment as shown in FIG. 30D. This makes micronization of devices possible.\nExamples of SAC are disclosed in Japanese Patent Laid-Open Publication No. 292323/1986, Japanese Patent Laid-Open Publication No. 106929/1992 and '94 Symp. VLSI Tech., Tech. Dig., pp. 99-100.\nNext, the conventional BLC will be explained in comparison with a structure which does not use BLC.\nIn the structure as shown in FIG. 31A in which a device isolation film 12 is formed on a semiconductor substrate 10, and an inter-layer insulation film 20 is formed, in opening a contact hole 22 near the device isolation film 12, the contact hole 22 and the device isolation film 12 must be spaced from each other so that the contact hole 22 is not positioned on the device isolation film 12 even when misalignment takes place.\nThat is, when the contact hole 22 is positioned on the device isolation film, the device isolation film 12 is adversely etched in opening the contact hole 22, and when a conducting film 24 is buried in the contact hole 22, a junction short circuit takes place between the conducting film 24 and the semiconductor substrate 10 (FIG. 31B).\nIn contrast to this, in the BLC as shown in FIG. 31C, an inter-layer insulation film 20 is formed of insulation films 16, 18 having etching selectivity different from each other. By forming the insulation film 16 in contact with the device isolation film 12 of a material having sufficient etching selectivity with respect to the device isolation film 12, the device isolation film 12 is not etched when a contact hole 22 is opened to expose the surface of the semiconductor substrate 10. A junction short circuit between the conducting film buried in the contact hole 22 and the semiconductor substrate 10 can be prevented.\nThus, in the BLC, even when the device isolation film 12 and the contact hole are overlapped, a junction short circuit can be prevented, and it is not necessary to consider an allowance for alignment between the device isolation film 12 and the contact hole 22. The contact hole 22 can be positioned as exemplified in FIG. 31D. This makes device micronization possible.\nSemiconductor devices having the above-described conventional BLC, however, have the following problems.\nIn etching the insulation film 16, it is preferred that wet etching is used to allow a good selectivity between the insulation film 16 and the device isolation film 12, but because wet etching is isotropic and adversely etches the insulation film 15 below the isolation film 18, a cavity is adversely formed. (FIG. 32A). The thus, formed cavity 30 cannot be covered by conventional sputtering and remains even after the conducting film 24 is deposited (FIG. 32B). When a plug 26 is formed in the next step of forming a contact, made by burying tungsten (W), then WF.sub.6, which is a source gas for tungsten, intrudes through the cavity 30 to cause substrate erosion called a worm hole. Junction spiking adversely occurs within a source/drain diffused layer 14 (FIG. 32C).\nIn addition, when using Al (aluminum) deposited by CVD as wiring material, in place of the above-described W plug, the Al and the semiconductor substrate directly contact each other in the cavity 30, and the Al and the semiconductor substrate react with each other by means of heat treatments in later steps, with a result that junction spiking takes place in the source/drain diffused layer 14 (FIG. 33A).\nIn addition, the above is true for a case in which Cu is used as a wiring material. Cu especially forms a deep energy level in the forbidden band when it is diffused in the semiconductor substrate which result much degrades the characteristics of the transistors. In addition, Cu is apt to diffuse in the silicon oxide film, and if the Cu reaches the gate oxide film 34, the Cu may add to leak current between the gate electrode 40 and the semiconductor substrate 10 (FIG. 33B).\nIn addition, in a semiconductor device shown in FIG. 34, which includes a contact plug 206 buried in an inter-layer insulation film 202 on a semiconductor substrate 200, it is often a case that when the BLC structure is used in opening a via-hole for connection to a wiring layer 210, an insulation film 220 is etched immediately onto the inter-layer insulation film 202 due to misalignment in opening the via-hole, etc., and a contact plug 208 is exposed in a cavity 224 formed when etching the etching stopper film 216, with a result that a contact plug 230 short-circuits with the contact plug 208.\nIn addition, when an etching stopper film 112 is removed without forming a cavity 124 using an anisotropic reactive ion etching, it is difficult to secure selectivity with respect to a ground film.\nThat is, in the structure of FIG. 35A, when the etching stopper film 112 in a wiring groove 118 is etched under conditions which secure a sufficient selectivity with respect to an inter-layer insulation film 104, a sufficient selectivity with respect to a contact plug 110, and the contact plug 110 is sometimes etched (FIG. 35B).\nOpposite to this, when the etching stopper film 112 is etched under conditions which secure a sufficient etching ration with respect to the contact plug 110, a sufficient selectivity with respect to the inter-layer insulation film 104, and the inter-layer insulation film 104 is sometimes etched (FIG. 35C).\nThus, in etching the etching stopper film 112, it is difficult to concurrently secure etching selectivity with respect to both the contact plug 110 and the inter-layer insulation film 104. Contact characteristics are degraded, which affects reliability of the semiconductor device.\nIn addition, when the BLC structure is used in forming a contact plug 144 on a wiring layer 122 buried in an inter-layer insulation film 114, the wiring layer 122 is exposed in a cavity 138 formed by retreating an etching stopper film 130. When a plug 142 is buried, a raw material gas of the plug 142 reacts with the wiring layer 122 to form a high-resistance reaction product 146, which sometimes deteriorates contact characteristics between the contact plug 144 and the wiring 122 (FIG. 36).\nThe inventors have found a new problem which had not been discovered, in studying semiconductor devices of the above-described conventional structures.\nThat is, in the SAC as exemplified in FIG. 37A in which a gate electrode 40 and a contact hole 22 overlap each other, and a step is formed in the contact hole, it has been found that when the contact hole 22 is opened in an inter-layer insulation film 20 formed of an insulation film 18 and an insulation film 16 which are formed of SiN film, the SiN film is more worn at the shoulder of the step in etching the insulation film 18. As a result, when the worn SiN film is removed using the conventional technique, even an insulation film 38 directly below the SiN film is adversely etched as shown in FIG. 37A by the dotted line, and a gate electrode 40 is adversely exposed.\nWhen phosphoric acid and fluorine radicals are used in the etching to thereby raise etching selectivity between the SiN film and the oxide film to suppress the above-described wearing of the oxide film, as shown in FIG. 37B, the etching horizontally advances in the insulation film 16, and a cavity 30 is adversely formed. Subsequently when a conducting film 24 is deposited, the conduction film 24 is not deposited in the cavity 30. In the next step of forming a contact by burying W, WF.sub.6 intrudes through the cavity 30 to cause substrate erosion called a worm hole. Junction breakdown adversely takes place near a source/drain diffused layer 14.\nAlso in the case in which a self-aligned silicide (salicide) is formed on the source/drain diffused layer 14, a semiconductor substrate 10 is not sufficiently covered with the silicide layer 44 at the edge of the device isolation film 12, and a worm hole is caused at the edge which results in junction breakdown (FIG. 38)."}
{"text": "Three-dimensional vertical NAND strings having one bit per cell are disclosed in an article by T. Endoh et al., titled “Novel Ultra High Density Memory With A Stacked-Surrounding Gate Transistor (S-SGT) Structured Cell”, IEDM Proc. (2001) 33-36.\nFluorine diffusion between word lines and insulating layers in an alternating stack of layers of a three-dimensional memory device can adversely impact reliability of the three-dimensional memory device by forming voids in the insulating layers, causing migration of materials, and forming electrical shorts. Fluorine diffusion needs to be suppressed in order to provide a reliable three-dimensional memory device."}
{"text": "The present invention relates to speech processing apparatus and methods. More particularly, the present invention relates to apparatus and methods for use in automatic speech recognition applications and research.\nSpeech, as it is perceived, can be thought of as being made up of segments or speech sounds. These are the phonetic elements, the phonemes, of a spoken language and they can be represented by a set of symbols, such as International Phonetic Association symbols.\nThese segments are linguistic units and have their bases in speech as it is perceived and spoken. All of the syllables and words of a language are made up of a relatively small number of phonetic elements. For example, in the case of English, textbooks in phonetics may list as few as 25 consonants and 12 vowels for a total of 37 phonemes. If the finer phonetic distinctions are included, then the list of distinguishable speech sounds or phones may lengthen to as high as 50 or 60.\nIt has been proposed that the phonemes of a spoken language can be understood in terms of a small set of distinctive features numbering about 12. These features have their bases in articulatory, perceptual, and linguistic analyses. A feature approach is often used in textbooks on phonetics as the phones and phonemes are described in terms of place of articulation and manner of articulation.\nThere are several theories of how the human listener processes an incoming acoustic waveform of speech and translates that waveform into a series of linguistic elements such as phonemes or words. The exact mechanisms and processes involved in the perception of speech are not yet fully understood. Finding simple and reliable acoustic-auditory correlates of the phones, phonemes and presumed features has proved elusive.\nResearch on speech perception has led to complicated, highly conditioned statements of relations between acoustic- auditory patterns and perception of phonemes, and even these statements are often of narrowly circumscribed generality. For example, the problem of how the listener can divide the acoustic input into segments relevant to linguistic perception is not understood. Even if a solution of this segmentation problem were available, the auditory-acoustic expression of a phoneme or feature seems to depend on the phonetic context, the particular talker, and the rate of speaking.\nAs a result of these problems there are several viable theories of speech perception. All of the current theories can be cast into a generic three-stage model, with the acoustic input undergoing three stages of processing in a bottom-up sequence. Stage 1 is an auditory-sensory analysis of the incoming acoustic waveform whereby representation of the signal is achieved in auditory-sensory terms. Stage 2 is an auditory- perceptual transformation whereby the spectral output of stage 1 is transformed into a perceptual form relevant to phonetic recognition. Here the spectral descriptions are transformed into dimensions more directly relevant to perception. For example, in various theories the perceptual form may be related to articulatory correlates of speech production or auditory features or pattern sequences. Finally, there is stage 3 in which the perceptual dimensions of stage 2 are transformed by a phonetic- linguistic transformation into strings of phonemes, syllables, or words. Stages 2 and 3 also are influenced by top-down processing wherein stored knowledge of language and events and recent inputs, including those from other senses as well as language, are brought into play.\nSome work in automatic speech recognition has involved a narrow-band spectral analysis performed on a time-windowed speech waveform. In one system described in \"Recognizing continuous speech remains an elusive goal\" by R. Reddy et al., IEEE Spectrum, Nov., 1983, pp. 84-87, incoming digitized signals are broken into centisecond slices and spectrally analyzed. Each slice is compared with a collection of sound prototypes and the prototype closest to each slice is entered into a sequence. The prototype sequence is then used to roughly categorize the initial sound of the word, which in turn is used to produce word hypotheses. Each word is then tested by creating a probability matrix and a cycle of operation repeats for the next word until an entire sentence is identified.\nDiphthongs, glides, and r-colored vowels are speech sounds that are all generically referred to as glides herein. Analysis of these sounds continues to pose difficult problems among the many faced in the field of automatic speech recognition. A paper which discusses some of these types of speech sounds is \"Transitions, Glides, and Diphthongs\" by I. Lehiste et al., J. Acoust. Soc. Am., Vol. 33, No. 3, March, 1961, pp. 268-277."}
{"text": "1. Field of the Invention\nThis invention relates to vegetable shelling apparatus and, in particular, to a vegetable shelling apparatus wherein the separated vegetables pass outwardly from a basket volume in a circumferential direction.\n2. Description of the Prior Art\nVegetables, such as peas or beans, when harvested, are encased within a typically dry or membranous outer protective covering known generally as the shell. In order to further process the vegetables it is necessary to separate the vegetable from its protective shell. However, care must be exercised to avoid damage to the vegetables during the husking, or shelling, operation.\nOf course, vegetables may be manually shelled, but such a process is burdensome, time-consuming and expensive. Therefore, the prior art provides several vegetable shelling devices, these devices being generally groupable into two main classes.\nThe first class of shelling apparatus provides an enclosed mesh or otherwise perforated basket, usually cylindrical in shape, into which the vegetables to be shelled are introduced. A suitable agitator is provided within the volume of the basket to break the shells and eject the vegetables therefrom. The openings in the mesh or perforations of the basket are usually sized to permit the shelled vegetables to pass therethrough, while the spent, empty shells are retained within the basket volume. The shelled vegetables pass radially outwardly (relative to the axis of the cylindrical basket volume) and may be collected in a tray, basket, or other receptacle disposed beneath the basket. Although the basket is usually rotated in a predetermined angular direction (no specified direction being optimum), in theory, the only movement required is the agitation of the husked vegetables, for once the shells are broken, the vegetables fall by their own weight through the openings in the perforated basket. Exemplary of such prior art devices are U.S. Pat. Nos. 496,206 issued to Placide, 421,244 issued to Chisholm et al., and 457,538 issued to Moulton.\nA variation on the basic theme is provided in U.S. Pat. No. 3,825,018, to Ferrars, which provides means for radially clearing the openings in the mesh or perforated basket to prevent the unwanted accumulation of spent shells from blocking the egress passages for the shelled vegetables. In the same vein is U.S. Pat. No. 3,041,486 issued to Farmer, Sr., which extends needle-like elements through the mesh or perforations to keep them clear and open. U.S. Pat. No. 3,087,499 issued to Carmichael, Jr., et al., describes a resilient mesh basket structure. U.S. Pat. No. 3,084,696, issued to Carmichael, Jr., relates to an agitator operatively connected to the basket which acts to prevent clogging of the mesh or perforated openings.\nA second class of prior art shelling apparatus is that exemplified by U.S. Pat. No. 3,238,950 issued to Taylor. This type of apparatus utilizes closely spaced, stacked roller elements which cause the vegetables to be ejected from their shells by the compression of the shells between the rollers. Such operation is analagous to the wringer operation commonly used on now-outdated washing machine apparatus. One disadvantage of the structure exemplified by the last-mentioned patent occurs due to the requirement that the vegetables be oriented so that the shells enter the rollers substantially perpendicular to their axes of rotation. It is also apparent that a portion of the vegetables being shelled may be squashed or otherwise destroyed with apparatus of the second general classification.\nIn view of the foregoing, it is advantageous to provide a vegetable shelling apparatus wherein the vegetables to be husked, or shelled, are inserted or deposited within a volume within a basket member such that agitation thereof causes the vegetables and the spent shells to collect in the radially outer regions of the interior of the basket. It is also of advantage to move strainer elements through the collected vegetables and their shells so that the vegetables pass circumferencially through the strainers and out of the basket while the shells are retained therein. Further advantage may be provided by the disposition of the strainer elements in radial gaps between the basket itself and the leading edges of arcuate segments projecting into the basket region. Still further advantage may be provided by rotation of the basket and the agitator at predetermined speeds, so as to maximize the ejection and separating of the vegetables from their shells."}
{"text": "A typical container refrigeration apparatus disclosed, for example, in Patent Document 1 cools air in a container used for marine transportation, for example.\nPlants (e.g., bananas and avocados) in such a container still breathe even after harvest by absorbing oxygen in the air and releasing carbon dioxide. These breathing plants consume nutrients and water stored in the plants themselves, and the plants significantly deteriorate as their respiration rate increases. Thus, oxygen concentration in the container is preferably low to the degree that the plants do not develop a breathing problem.\nPatent Document 1 discloses an internal environment control system to reduce oxygen concentration of air inside a container by generating nitrogen-enriched air higher in nitrogen concentration than atmosphere and supplying this nitrogen-enriched air into the container. When the internal environment control system reduces the oxygen concentration of the inside air, a respiration rate of plants in the container decreases, curbing deterioration of the plants."}
{"text": "With development of science and technology, an electronic product with a biometric identification function gradually enters people's life and work. Due to uniqueness and invariability, a fingerprint may be used for personal identification. Therefore, a fingerprint identification technology has attracted much attention."}
{"text": "Gear transmission speed change devices are capable of high torque but are complex to manufacture, exhibit backlash, gear noise, and typically require multiple stages to achieve high reduction ratios.\nTraction drive speed change systems offer certain advantages but typically produce less torque then a geared speed change device of the same size.\nGenerally speaking, with traction drive systems it is desirable, for many applications, to provide pure rolling contact between torque transmitting members for low friction, high efficiency, long service life, and increased traction. Furthermore, is desirable for the contact forces between the torque transmitting members to increase automatically as torque increases, so high loads are only generated between traction components when necessary to allow increased torque output. This would provide the benefit of increased service life and efficiency by reducing wear surface loading when the actuator is not transmitting high torque.\nOther desirable characteristics of a rotary actuator are zero backlash, backdriveability, low vibration, non-cogging output, high rigidity, and quiet operation. High torque capacity for size and weight are also desirable, as are a wide range of speed change ratio possibilities (including very high speed change ratios) high input speed capability (to allow the use of low torque input drive systems) low input inertia, and a relatively large center thru-hole for interned wiring (or possibly to allow an integrated electric w other type of rotary drive motor within (or partially within) the inner diameter (ID) of the actuator."}
{"text": "Field\nThis disclosure relates generally to implants, and, more specifically, to hydrogel joint implants and various tools, devices, systems, and methods related thereto.\nDescription of Related Art\nImplants can be used to replace deteriorated or otherwise damaged cartilage within a joint. Such devices can be used to treat osteoarthritis, rheumatoid arthritis, other inflammatory diseases, generalized joint pain, joints damaged in an accident, while damaged participating in athletics, joints damaged due to repetitive use, and/or other joint diseases."}
{"text": "1. Field of the Invention\nThis invention relates to a zoom lens, and more particularly, to a zoom lens which is suitable for a still-picture camera, which includes a wide-angle range, and which has a high varifocal ratio.\n2. Description of the Prior Art\nRecently, zoom lenses have generally been used as interchangeable lenses for single-lens reflex cameras. This tendency is a result of the high performance, size reduction and low price of zoom lenses. Zoom lenses which have greater varifocal ratios, and which are even more compact, hereafter also will be required. In the field of compact cameras, where lenses are built into the camera, zoom lenses which have greater varifocal ratios and which are more compact also will be required.\nFurthermore, in accordance with the popularity of autofocus cameras, optical systems suitable for an autofocusing operation also are required. Hence, it has become required that the amount of movement, weight, diameter and the like of a focus-moving lens have smaller values.\nZoom lenses having high varifocal ratios mainly are used for single-lens reflex cameras and have been disclosed, for example, in Japanese Patent Publication No. 58-33531 (1983) and Japanese Patent Laid-open Application (Kokai) Nos. 59-198214 (1984) and 60-39613 (1985). So-called rear-focusing zoom lenses, in which focusing is performed by lens groups other than the front lens group also have been disclosed, for example, in Japanese Patent Laid-open Application (Kokai) Nos. 58-136012 (1983); 58-143312 (1983); 61-138228 (1986); 63-205628 (1988); 63-205629 (1988); 63-266414 (1988); 63-266415 (1988) and 1-154014 (1989).\nThese zoom lenses have the configuration of five groups having positive, negative, positive, negative and positive refractive powers, arranged from the object side, in order to efficiently perform a varifocal operation and to provide a small size. However, these efficiencies are still insufficient. While the diameter of a front lens can be reduced by adopting a rear-focusing configuration, such a configuration is so complicated that the response speed of an autofocusing operation becomes insufficient."}
{"text": "1. Field of the Disclosure\nThe invention is directed to a process, system and device to provide service technicians with augmented reality service information and, in particular, augmented reality service information that provides the internal and external configurations, schematics and details of components of vehicles and the like.\n2. Related Art\nCurrently service technicians must use hard copy diagrams or computer-based diagrams in order to schematically view various internal aspects of a vehicle for repair purposes. Currently there is no solution to provide the service technician with the ability to move a computer with respect to a vehicle and have the schematics and diagrams change in conjunction with the location of the computer with respect to the vehicle. Moreover, there is no solution to provide the service technician with the ability to move a computer with respect to a vehicle and have the schematics and diagrams change in conjunction with the location of the computer with respect to the vehicle such that they are superimposed over the vehicle image as shown in the computer.\nAccordingly, there is a need for a computer that can provide schematics and diagrams in response to movement of computer with respect to a vehicle to improve the repair process."}
{"text": "The subject matter disclosed herein generally relates to the field of control systems, and more particularly to a method and apparatus for a configurable sever system.\nCritical embedded controllers typically include sever logic that disables controller outputs when normal processing functions of the controller have malfunctioned. Sever logic varies between controller designs and is typically customized in a fixed format for a desired end use or application."}
{"text": "The present invention is related to the field of optical switches and, more particularly, to optical switches in which the switching function is effected by electromechanical devices.\nIn a 1.times.N optical switch, a single optical fiber forms an optical path with an optical fiber selected from N optical fibers. Thus a signal on the single optical fiber is transmitted to the selected fiber. Similarly a signal on the selected fiber is passed to the single optical fiber. A switch permits the optical connection to be changed from the selected fiber to another of the N fibers.\nIn electromechanical 1.times.N optical switches, the switching operation is effected by the movement of an end of the single optical fiber relative to the ends of the N optical fibers. To form an optical path, the end of the single optical fiber is aligned with the end of the selected fiber from the N fibers. In such systems alignment of the fiber ends is critical and the much effort is required with the electromechanical drive devices, such as motors and their movements, to ensure precise and reproducible alignment of the single fiber with each of the N optical fibers. Nonetheless, performance of such switches has remained poor with use restricted to multimode optical fibers. In present day optical networks it is often desirable to use single mode optical fibers.\nThe present invention presents an electromechanical 1.times.N optical switch with high optical performance. Moreover, with the present invention single mode optical fibers may be used."}
{"text": "Contaminated water exists in ever increasing quantities. Some of these waters contain low levels of organic compounds such as 2,4-dinitrotoluene, 1,4-dioxane, ethylene glycol, trinitrotoluene, phenols, anilines and pentachlorophenol. These compounds are susceptible to oxidation and/or photooxidaton processes. By such reactions, these organic compounds can be oxidized to carbon dioxide and water and, in the case of nitro or chloro compounds, NO.sub.3.sup.- and Cl.sup.-.\nAccordingly, when oxidizable compounds are present in relatively low concentrations, less than 10,000 ppm, the water may be treated by the addition of one or more oxidizing agents accompanied, in some instances, with irradiation. Oxidizing agents which result in the production of hydroxyl and hydroperoxy radicals are typically employed. These radicals are very strong oxidizing agents and, in addition, they indiscriminately oxidize all organic compounds. Typical sources of such oxidizing radicals are ozone and hydrogen peroxide. The irradiation is normally conducted by using ultraviolet light having a wave length of from about 200 to about 300 nm. The choice of oxidizing agent and the use of ultraviolet light will depend upon various factors including the concentration of oxidizable contaminants in the water stream, the required amount of oxidizing radical and the pH of the water.\nThe oxidation/photooxidation treatment may be conducted to any stage of completion which is required. If the only contaminants present in the water are oxidizable contaminants, the water may be treated so that it meets environmental requirements for use or disposal. Alternatively, the water may be recycled or, if other non-oxidizable contaminants are present, subjected to further treatment.\nTypically, waters can be held in holding tanks prior to the oxidation/photooxidation treatment. In such cases, the water is exposed to carbon dioxide. As a result of this exposure, carbonate and bicarbonate ions are formed. These ions compete for the oxidizing hydroxyl radicals. As a result of this competition, the amount of oxidizing agents to be added to the solution may have to be increased if the treatment rate is not to be greatly diminished. However, if ozone is being used, then it may not be possible to increase the flow rate of ozone to the oxidation reaction without increasing the size of the ozone plant. There is therefore a need for an efficient method to control the competition for the oxidizing radicals by the carbonate and bicarbonate ions."}
{"text": "1. Field of the Invention\nBisphosphonates--synthetic compounds containing two phosphonate groups bound to a carbon and two additional groups R.sub.1 and R.sub.2, respectively (see formula A herebelow)--bind strongly to calcium crystals, inhibit their growth, suppress bone resorption and are used in the treatment of a variety of disorders of calcium and bone metabolism. It is generally considered that the \"bone hook\" (P-C-P with R.sub.1) is responsible for the binding of these molecules to mineralized matrices, while, in turn, the R.sub.2 group is primarily responsible for its effect on bone resorption. ##STR2##\n2. Description of the Related Art\nU.S. Pat. No. 4,054,598 discloses a method with sequestering agents, especially for alkaline earth metal ions, having the formula 1-hydroxy-3-amino-alkane-1,1-bisphosphonic acids, useful for the treatment of disturbances of calcium or phosphate metabolism characterized by abnormal deposition of difficultly soluble calcium salts or the abnormal dissolution of hard tissues causing losses of hard bone substance, which cannot be replaced or only by incompletely crystallized tissues, such as Paget's disease, lithiasis, arthritis and others. Spanish Patent No. P910088 discloses pharmaceutical liposome preparations containing bisphosphonates as active compounds. The use of olpadronate as an anabolic agent for the treatment of osteoporosis and other bone metabolism disorders is claimed in the Applicant's earlier PCT application PCT/EP95/05142.\nA number of 1,3-diaminoalkane-1,1-bisphosphonic acids including, amongst others, 1-amino-3-(N,N-dimethylamino)-propyliden-1,1-bisphosphonic acid are disclosed in DE-A1-25 34 390. Several applications (medical and non-medical), unrelated to the claimed methods of use of the present application, are mentioned, without giving any specific advantages for 1-amino-3-(N,N-dimethylamino)-propyliden-1,1-bisphosphonic acid.\nDE-A1-18 13 659 is directed to compositions for the inhibision of deposition and mobilisation of calcium phosphate in tissues with an effective, but non-toxic content of various polyphosphonates. However, the specific compound of the present application is not mentioned.\nIn Eur. J. Clin. Invest. 1970, Vol. 1 (1) pp. 12-18 (Fleisch et al.) the inhibitory effect of phosphonates on the formation of calcium phosphate crystals in vitro and an aortic and kidney calcification in vivo is described. However, the specific compound of the present application is not mentioned.\nDE-C1-10 02 355 is directed to processes for the production of various bisphosphonates. However, the specific compound of the present application is not mentioned. Moreover, no specific applications for the described compounds are given.\nZ. Anorg. Allg. Chem. 389, pp. 119-128 (1972) (Ploger et al.) is directed to different processes for the production of various bisphosphonates, however without mentioning the specific compound of the present application or any specific application of the described compounds at all.\nCA-A1-21 20 538 discloses the use of bisphosphonic acid derivatives for promoting bone repair. Although the general formula would also cover the specific compound of the present application, this compound or any related advantages are not mentioned in the description.\nIt is the object of the invention to provide for specific new medical applications of 1-amino-3-(N,N-dimethylamino)-propyliden-1,1-bisphosphonic acid or its monosodium or other pharmaceutically acceptable salt."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for producing an L-amino acid such as L-cysteine. More precisely, the present invention relates to a bacterium suitable for production of an L-amino acid and a method for producing an L-amino acid utilizing such a bacterium. L-Amino acids are used in various fields, for example, as seasonings, food additives, animal feed additives, chemicals, drugs, and so forth.\n2. Background Art\nL-Amino acids are industrially produced by fermentation using microorganisms belonging to the genus Brevibacterium, Corynebacterium, Escherichia, or the like. In such methods, strains are used which are isolated from nature or artificial variants of such strains. Furthermore, microorganism strains can be used which are modified by recombinant DNA techniques so that the activity of a basic L-amino acid biosynthesis enzyme is increased, and so forth (EP 0643135 B, EP 0733712 B, EP 1477565 A, EP 0796912 A, EP 0837134 A, WO01/53459, EP 1170376 A, WO2005/010175, WO96/17930, WO2006/013807).\nFurthermore, L-cysteine, for example, is obtained by extraction from substances containing keratin, such as hair, horns, and feathers. L-cysteine can also be obtained by converting the precursor DL-2-aminothiazoline-4-carboxylic acid using a microbial enzyme (Sano, K. et al., Appl. Environ. Microibol., 34, 806-810 (1977)).\nFurthermore, L-cysteine has also been produced by fermentation utilizing a bacterium. For example, a method has been disclosed for producing L-cysteine using an Escherichia bacterium with a suppressed L-cysteine decomposition system and a serine acetyltransferase (EC 2.3.1.30, henceforth also referred to as “SAT”) with attenuated feedback inhibition by L-cysteine is attenuated (Japanese Patent Laid-open (Kokai) No. 11-155571). Furthermore, coryneform bacteria or Escherichia bacteria attenuated or deleted activity of cystathionine-β-lyase (Japanese Patent Laid-open (Kokai) No. 11-155571), tryptophanase (Japanese Patent Laid-open No. 2003-169668), or O-acetylserine sulfhydrylase B (Japanese Patent Laid-open No. 2005-245311) are examples of bacteria with an enhanced ability to produce L-cysteine by suppressing the decomposition of L-cysteine. A method for producing L-cysteine by using a bacterium in which L-cysteine metabolism is decontrolled by using a DNA sequence coding for SAT that has a specific mutation for attenuating feedback inhibition by L-cysteine is also known (National Publication of Translated Version in Japan (Kohyo) No. 2000-504926).\nFurthermore, it is known that the ydeD gene which encodes the YdeD protein (Dabler et al., Mol. Microbiol., 36, 1101-1112 (2000)), and the yfiK gene which encodes the YfiK protein (Japanese Patent Laid-open No. 2004-49237), participate in secretion of the metabolic products of the cysteine pathway. Techniques are also known for enhancing the L-cysteine-producing ability by increasing expression of the mar-locus, acr-locus, cmr-locus, mex-gene, bmr-gene, or qacA-gene, which encode proteins responsible for secreting toxic substances from cells (U.S. Pat. No. 5,972,663), or emrAB, emrKY, yojIH, acrEF, bcr, or cusA genes (Japanese Patent Laid-open No. 2005-287333).\nMoreover, a method has been reported for producing L-cysteine using a bacterium which overexpresses a gene coding for a protein responsible for releasing an antibiotic or a substance toxic to a bacterium directly from a cell (Japanese Patent Laid-open No. 11-56381).\nMoreover, it is known that production of L-amino acids by bacteria, not only L-cysteine, can be improved by increasing expression of an L-amino acid secretion protein. It is known that, in coryneform bacteria, L-lysine production is improved by increasing expression of an L-lysine secretion carrier, named LysE (Japanese Patent Laid-open No. 2000-189180). Furthermore, for Escherichia bacteria, several membrane proteins are known to be amino acid secretion carriers. For example, it has been reported that, by increasing the copy number of a gene named rhtB, resistance to a high concentration of L-homoserine, L-threonine, L-alanine, L-valine, and L-isoleucine is improved, and the product of this gene is a secretion carrier of these L-amino acids (Japanese Patent Laid-open No. 2000-189180).\nThe yeaS gene belongs to the family of the aforementioned rhtB gene, the product of which is presumed to be a membrane protein, and it has been reported that, by increasing expression of this gene, resistance to a high concentration of L-threonine, L-homoserine, L-lysine, L-glutamic acid, L-histidine, L-proline, and α-aminobutyric acid is improved as compared to a control strain, and by increasing expression of this gene in an L-amino acid-producing bacterium, production of L-valine, L-isoleucine, L-alanine, L-proline, and L-histidine is improved (EP 1013765 A1).\nAs described above, the effect of increasing yeaS gene expression on the production of various amino acids has been investigated (Kutukova et al., FEBS Lett., 579, 4629-4634 (2005)). However, mutations of the yeaS gene which improve L-amino acid productivity have not been previously reported."}
{"text": "The combustion of coal in a boiler, as in a pulverized coal-fired electric power generating plant, produces flyash. The composition of the flyash varies depending, for example, on the composition of the coal and the combustion conditions. Generally, flyash is a fine, solid, noncombustible mineral residue, which is distinct from bottom ash, cinders, or slag. Flyash can have widely varying particle size, density, shape, porosity, internal structure, and surface chemistry. It is typically composed of oxidized silicon, aluminum, calcium, iron, titanium, magnesium, sodium, potassium, sulfur, etc.\nThe sources of flyash from coal can generally be classified into two categories: mineral inclusions, i.e., extraneous minerals; and, organically associated inorganic elements (\"OAI's\" or inherent minerals). Inherent minerals are the components of the coal, such as sulfur, sodium, calcium, and potassium, which are not present as mineral inclusions in the coal matrix, but are actually associated with the chemical structure of the complex hydrocarbons which make up the coal's combustible component. The mineral inclusions are the solid, generally crystalline, compounds that are found in salt, rock, clay, and iron pyrites, for example.\nThe formation of flyash during coal combustion generally depends upon the transformation of minerals during the pyrolytic process of combustion, and the release of inorganic elements on an atomic or near atomic scale from the hydrocarbon matrix that comprises the structure of the coal itself. The inorganic elements released from the organically associated inorganics form a \"fume,\" i.e., a suspension of particles in a gas, with an average particle size of about 1 micron or less. In some instances, the minerals which form the fume will be such as to exist in the vapor state, at least when the fume is the hottest. This may be the case, for example, for sodium and potassium oxides. Typically, the fume is composed of oxides of such elements as sodium, calcium, potassium, and magnesium.\nMineral associated flyash, i.e., ash formed from the mineral inclusions, commonly exists in a fairly wide range of particle sizes. Generally, however, it is most often between about 1 micron and about 100 microns in size (diameter). That is, the particle size distribution of the flyash formed from the mineral inclusions is typically such that the bulk of it, by weight, is of particles about 2 to 70 microns in diameter. In flyash, such materials are often generated from coal as glassy cenospheres.\nDisposition of flyash from coal-burning installations such as power generating plants is an increasingly difficult problem. Strict environmental restrictions pertaining to air quality standards and the handling and final placement of flyash have combined to make flyash a source of escalating processing costs and environmental concerns common to nearly all coal-burning plants. To meet the environmental standards, flyash is generally removed from the exiting coal combustion off-gases by such arrangements as scrubbers or baghouses. In a typical example, the gas is fed through a shower of water, such as droplets in a venturi scrubber (or aqueous scrubber). The flyash is collected by the water as the gas passes therethrough. The gas is thereby cleansed and the particulate matter in the water is collected or settled in a pond.\nStack \"opacity\" is a government regulated flyash emission parameter. It generally concerns definition of the \"clarity\" of stack emission; i.e., percent transmission through a volume of stack emissions. The greater the opacity, the more contaminated the emissions. Extreme stack opacity values can limit the types of coal and/or amount of power that can be produced at a generating unit. That is, certain types of coal cannot be burned without extremely efficient scrubber systems or reduced power output because they generate a large amount of particulate matter, which contributes to opacity. Therefore, some coal-burning facilities are limited in the types of coal that can be burned in order to meet particulate emission standards.\nWhat has been needed is still further systems and methods for reduction in the amount of flyash emissions from combustion processes. Such systems and methods would allow a wider range of coals to be burned without penalty, resulting in a more aggressive coal fuel purchasing strategy, and reduced cost of electricity production. The particulate emissions can be reduced; and, the power plant can regain a greater total power output (within opacity limits), if it was \"opacity limited.\" Other advantages may result, such as reduced sulfur and/or flyash output from the plant as a result of the properties of the new coal used.\nIn addition to flyash emission problems, coal-burning facilities are faced with ash fouling problems. This is because coal-burning facilities have become more efficient by increasing the temperature of the steam produced in the boiler. Boilers, and the tubing (heat exchange surfaces) in the boilers, have also been improved so as not to be the limiting factor in obtaining these high temperatures. However, if the boiler tubes (or heat exchange surfaces) are so hot that they exceed the fluxing temperatures of the flyash which is being transported through the tubes along with the combustion off-gases, the flyash can adhere to the tubes. The flyash deposits can then build up in the tubes and interfere with the movement of off-gases and the rate of steam production. This detrimentally effects the efficiency and capacity of the boiler.\nCertain types of coal that produce a relatively low amount of flyash upon combustion can be burned, with concomitant reduction in this ash buildup, i.e., ash fouling, problem. However, this is not always economically efficient. It has also been suggested that vermiculite can be added to the gases and flyash produced during a combustion process. This method, however, does not prevent the formation of flyash deposits. The vermiculite actually combines with the flyash to form ash deposits. Although these deposits are easier to remove than pure ash deposits due to the ability of the vermiculite to expand when exposed to elevated temperatures, they must still be removed by the application of jets of steam or soot blowers. It is, therefore, generally desirable to develop a system and method that reduces the amount of flyash buildup in the boiler system."}
{"text": "For example, an image forming apparatus such as a color copying machine or a printer includes an exposing device. The exposing device includes a laser diode as a light source, a polygon mirror, a deflection lens (an fθ lens), and a reflection mirror. The exposing device may be referred to as a laser scanning unit or an optical scanning device. For example, as disclosed in U.S. Pat. No. 6,700,687 B1, a light source, a polygon mirror, a deflection lens, a reflection mirror, and the like are housed in a casing. A laser beam emitted from the light source is transmitted through the deflection lens via the polygon mirror and reflected by the reflection mirror to travel to a photoconductive drum.\nThe reflection mirror has a body square-bar-shaped in section and a reflection layer formed by depositing aluminum on the body. The mirror has a first surface on which the reflection layer is formed, second and third surfaces perpendicular to the first surface, and a fourth surface on the opposite side of the reflection layer. One end in an axial direction (a longitudinal direction) of the mirror is supported by one mirror supporting section provided in the casing. The other end in the axial direction of the mirror is supported by the other mirror supporting section provided in the casing. The exposing device of this type needs to reflect the laser beam at a predetermined angle with respect to the photoconductive drum. Therefore, the exposing device includes an adjusting mechanism for adjusting an angle of the mirror.\nFor example, the exposing device in the past, one mirror supporting section for supporting one end of the mirror has a positioning protrusion that comes into contact with the first surface of the mirror and an adjusting screw provided beside the positioning protrusion. A distal end of the adjusting screw is in contact with the first surface. A first flat spring is in contact with the fourth surface of the mirror with elastic energy stored therein. The first flat spring urges an end of the mirror toward the positioning protrusion. A positioning section that comes into contact with the second surface of the mirror is provided in one mirror supporting section. A second flat spring is in contact with the third surface of the mirror with elastic energy stored therein. The second flat spring urges one end of the mirror toward the positioning section. The other end of the mirror is elastically supported by a flat spring and a positioning protrusion provided in the other mirror supporting section.\nTo adjust an angle of the mirror, a projecting amount of the adjusting screw with respect to the mirror is changed by rotating the adjusting screw. For example, when the projecting amount of the adjusting screw is increased, while bending of the first flat spring increases, the mirror rotates in a first direction with the respective positioning protrusions as fulcrums. When the projecting amount of the adjusting screw is decreased, the mirror rotates in a second direction with the positioning protrusions as fulcrums with the repulsion of the first flat spring.\nIn the adjusting mechanism in the past, it is difficult to keep the balance of loads between the flat springs and the adjusting screw. For example, if the repulsion of the flat spring is too strong, when the projecting amount of the adjusting screw is increased, force larger than necessary force is applied to the mirror. This causes breakage of the mirror and shift of a position of the mirror from a predetermined position. Conversely, if the repulsion of the flat spring is too weak, when the projecting amount of the adjusting screw is reduced, the flat spring cannot push back the mirror. Therefore, the mirror cannot be rotated in a desired direction. Further, if the repulsion of the flat spring is weak, for example, a position of the mirror is not stabilized and the mirror moves when impact is applied to the mirror. In particular, when “permanent set in fatigue” due to aged deterioration of the flat spring is taken into account, it is difficult to properly keep the balance of loads between the repulsion of the flat spring and the adjusting screw.\nBoth the ends of the mirror are supported by a supporting section of the casing of the exposing device. Besides the mirror, a light source, a polygon mirror, a polygon motors a motor driver, and the like are housed in the casing. These act as heat generation sources during the operation of the exposing device. The casing of the exposing device may be affected by heat generated from peripheral apparatuses.\nA casing of synthetic resin can be easily molded in a desired shape, light in weight, and low in cost compared with a casing of metal. However, the casing of synthetic resin has small thermal conductivity compared with the casing of metal. Therefore, a thermal expansion amount near a mirror supporting section close to a heat source and a thermal expansion amount near a mirror supporting section distant from the heat source may be different. This causes an adverse effect on a quality of an image in an exposing device including plural mirrors that are used in a color copying machine.\nFor example, in a mirror supporting section in a position close to a light source, since a thermal expansion amount is large, positional deviation of a mirror is large. An amount of shift of a beam reflected by the mirror is large. On the other hand, in a mirror supporting section in a position distant from the light source, since a thermal expansion amount is small, positional deviation of a mirror is small. An amount of shift of a beam reflected by the mirror is small. As a result, the amount of shift of the beams is different for each of the mirrors. This causes color drift on a photoconductive member. In that case, adjustment work is necessary for measures against the color drift. Therefore, for example, printing is unfavorably interrupted.\nIf the casing is integrally molded by using a material having high thermal conductivity such as an aluminum alloy, the thermal expansion difference between the mirror supporting sections can be reduced. However, the casing of metal is high in cost compared with the casing of synthetic resin. As another measure, it is also conceivable to fix a press-formed sheet metal member to the casing of synthetic resin. This sheet metal member is formed in a predetermined shape by pressing a metal plate. A mirror supporting member is provided in the sheet metal member. However, in the sheet metal member in which a bent section and the like are formed by pressing, it is difficult to highly accurately form the mirror supporting section in terms of a position and a shape."}
{"text": "As functional complexity increases, or business rules change, the backbone process of computer programs deprecate, made useless by dependencies between decision-logic and task-logic, or made useless by alterations to the data signature that allows interactions between programs.\nThe current state of the art is starting from scratch with a completely new set of programs, leaving the defects and deficiencies of legacy programs behind, while creating a completely new set of the same. Typically, the protocols for program organization are centered on the naming of programs and program properties.\nPrograms are generally written to account for “logical steps” of a given process, which gives rise to the tendency to write lengthy programs consisting of decision-logic mixed with task-logic. Alterations to the expected programming process may mandate rewriting entire sections of code.\nAn alteration to a single business rules may indicate the need to rewrite one ore more logical processes, to account for the correct implementation of a process series, or sequence, of decisions and related tasks. The convolution of programming logical steps is compounded by the need for cross-program communications, a need usually handled by a series of parameters in the originating program handled by a series of matching parameters in the receiving program. Any change to the expected programming process may also affect the content, sequence, and naming of parameters used in program communication.\nAdding new functionality to an existing application is a concept referred to as layering programs. This means new programs need to be inserted at specific points in the process, and some existing programs may need to be updated. The way to implement a new “layer” is to begin the process of rewriting logical steps represented by one or more program modules, accounting for alterations to parameters between programs, and accounting for alterations to the sequence of program execution.\nThe long term benefits of the traditional approach to computer programming fail to materialize when attempting to integrate legacy program applications with newly developed modules."}
{"text": "1. Field of the Invention\nThe present invention relates to an antenna duplexer mainly used for a high-frequency circuit or the like of a radio system to share an antenna by a transmitter and a receiver.\n2. Description of the Prior Art\nBecause mobile communication has recently advanced, an antenna duplexer is used for a lot of portable telephones and automobile telephones. An example of the above conventional antenna duplexer is described below while referring to the accompanying drawings.\nFIG. 13 shows an exploded perspective view of a conventional antenna duplexer. In FIG. 13, symbols 1301 to 1306 denote dielectric coaxial resonators, 1307 denotes a coupling board, 1308 denotes a metallic case, 1309 denotes a metallic cover, 1310 to 1312 denote series capacitors, 1313 and 1314 denote inductors, 1315 to 1318 denote coupling capacitors, 1321 to 1326 denote coupling pins, 1331 denotes a transmission (hereafter TX) terminal, 1332 denotes an antenna terminal, 1333 denotes a receiving (hereafter RX) terminal, and 1341 to 1347 denote electrode patterns formed on the coupling board 1307.\nThe dielectric coaxial resonators 1301, 1302, and 1303 and the series capacitors 1310, 1311, and 1312, and inductors 1313 and 1314 constitute a TX band rejection filter. Moreover, the dielectric coaxial resonators 1304, 1305, and 1306 and coupling capacitors 1315, 1316, 1317, and 1318 constitute a RX band pass filter.\nOne end of a TX filter is connected to the TX terminal 1331 electrically connected with a transmitter and the other end of the TX filter is connected with one end of a RX filter and also connected to the antenna terminal 1332 electrically connected to an antenna. The other end of the RX filter is connected to the RX terminal 1333 electrically connected with a receiver.\nOperations of the antenna duplexer constituted as described above are described below.\nFirst, the TX band rejection filter shows a small insertion loss for a TX signal in a TX frequency band and makes it possible to transfer the TX signal from the TX terminal 1331 to the antenna terminal 1332 almost without attenuating the TX signal. Moreover, the TX band rejection filter shows an operation that RX signals input through the antenna terminal 1332 return to the RX band pass filter because the TX band rejection filter shows a large insertion loss for the RX signals in a RX frequency band and most input signals in the RX frequency band are reflected.\nHowever, the RX band pass filter shows a small insertion loss for a RX signal in a RX frequency band and makes it possible to transfer the RX signal from the antenna terminal 1332 to the RX terminal 1333 almost without attenuating the RX signal. Moreover, the RX band pass filter shows an operation that TX signals coming through a TX filter are sent to the antenna terminal 1332 because the RX band pass filter shows a large insertion loss for TX signals in a TX frequency band and most input signals in the TX frequency band are reflected.\nAn antenna duplexer used for a high-frequency band of mobile communication has wide band characteristics. Therefore, to secure a necessary attenuation value in a wide band, it is necessary to further increase the number of stages of cascaded dielectric coaxial resonators.\nIn the case of the above structure, however, when the number of stages of resonators is increased to increase the attenuation value, the loss in a signal pass band width increases. To avoid the bad effect, it is considered to increase the unloaded Q of a dielectric coaxial resonator. However, to increase the unloaded Q, it is necessary to increase the size of the dielectric coaxial resonator. This is reciprocal to the recent antenna-duplexer downsizing trend."}
{"text": "1. Field of the Invention\nThe present invention relates to a method to measure, install and fix an object, particular a method to measure, install and fix a toilet seat.\n2. Description of the Prior Art\nCurrently, an ordinary toilet seat or a computerized toilet seat to be installed on a flush toilet should make a holder fixed at two mounting holes of the flush toilet with bolts, washers and nuts by a worker using a wrench or a screw driver in an inconvenient or time-consuming course of assembling or disassembling the toilet seat or the computerized toilet seat that may be unseparated in the case of bolts or nuts rusted.\nIn addition, a holder of any toilet seat or computerized toilet seat being available in the market and featuring its particular specifications should be lifted gently to check secure fastening after the toilet seat or the computerized toilet seat was installed on a flush toilet. In order to provide a mounting base of the computerized toilet seat available to any flush toilet made by an unspecified manufacturer, some tolerances have to be designed on a holder's mounting holes for flexible adjustment within certain extents. However, it is possible to open more mounting holes by a user in the case of two mounting holes' positions beyond the tolerances; otherwise, the toilet seat fails to be assembled."}
{"text": "The present invention relates to an access method and apparatus for rapidly retrieving desired information from a high density record medium which stores information signals thereon, and more particularly to an access method and apparatus which rapidly and optically retrieves desired information from a rotating optical disk on which information signals are recorded concentrically or spirally at a high density.\nThe optical disk has a number of concentric or spiral tracks formed thereon at a predetermined pitch, and a number of sectors for delimiting data are formed for each track. In order to record external information at any position or reproduce or erase information recorded at any position, one of the tracks on the disk is sought and one sector on the track is then sought. That is, a macro-seek control (i.e., speed control and position control) in which a light spot is rapidly moved to a neighborhood of a target track, a tracking servo in which the light spot is maintained at a center of the track, and a micro-seek control (for example, jump control) in which a deviation from the target track is corrected are necessary. The access operation in such an optical disk storage is disclosed in Japanese Unexamined Patent Publication Nos. 58-91536 and 58-169370 (corresponding to U.S. patent application Ser. No. 06/443,399, continuation application Ser. No. 06/736,125, now U.S. Pat. No. 4,607,358 and EPC Patent Application No. 82110907.1). In the optical disk storage, the positioning of the light spot is controlled by cooperation of a coarse actuator, such as a linear motor and a fine actuator, such as a galvanomirror. The coarse positioning of the optical head is carried out by the course actuator and an external scale as a position detector, and the fine positioning to the track is carried out by the fine actuator by detecting an address signal recorded on the track. A scale pitch of the external pitch is a multiple of the track pitch.\nThe macro-seek operation for positioning the optical head to an approximate target position by using the external coarse scale usually effects speed control to reduce an access time. However, since the macro-positioning is carried out by using the external scale which is independent from the track eccentricity, the deviation from the target track varies significantly due to the eccentricity. As a result, the distance of movement of the light spot in the next micro-seek control (for example, jump control) is large and the access time increases. In order to control the speed, a means for detecting the speed of the head is necessary. However, when a separate detector (such as a tachometer) for detecting the speed is mounted, the weight of the head increases and an oscillation at a low frequency may be generated depending on the length of the tachometer. In addition, the cost increases. Accordingly, it is prederable to detect the speed by using the output of the external linear scale. A pulse frequency of a zero-crossing detection signal is proportional to the speed of the head. Accordingly, the zero-crossing detection signal is frequency-voltage converted so that the speed is detected based on a voltage. However, the frequency-voltage conversion does not provide an accurate conversion output when the frequency is low (head speed is low), because the pulse interval of the zero-crossing detection signal increases and the speed change in the interval cannot be detected. Accordingly, when the speed is controlled by using the output of the frequency-voltage conversion, the control is unstable when the speed is low and the system may overrun."}
{"text": "There are numerous enzymes which have been identified as capable of catalyzing various chemical reactions. Similarly, it has been discovered that antibodies can be elicited to catalyze a variety of chemical reactions (U.S. application Ser. No. 674,253). It is well known that antibodies and enzymes share a fundamental similarity in that both are specialized proteins that bind to other molecules. However, there are important differences between antibodies and enzymes.\nAntibodies typically bind to a molecule or antigen so that the antigen is marked as foreign to the organism that produced the antibody. The binding of the antibody to the antigen enables the antigen to be removed from the organism. Enzymes are biological catalysts which bind a molecule in such a way that the activation energy of a reaction involving a molecule or substrate is lowered, thereby increasing the rate of the reaction.\nLinus Pauling discovered there are two types of interactions between proteins and the molecules that bind them. Antibodies bind molecules in their ground state while enzymes bind molecules in higher energy states.\nPauling attempted to explain the mechanism of enzyme catalysis based upon such binding. During the course of the chemical reaction, the reactants undergo one or more transitions through intermediate structures or transition states which are energically less favorable than either the reactant or the product. The hydrolysis reaction of a peptide linkage or an ester bond in an aqueous medium passes through a tetrahedral carbon transition state, as depicted in FIGS. 1 (peptide) and 2 (ester). In the transition state, a tetrahedral carbon atom is bonded to: a carbon atom of the acid portion of the peptide linkage or ester bond; two oxygen atoms, one corresponding to the carbonyl group and the other corresponding to a hydroxyl ion or water molecule of the medium; and either the oxygen atom of the alcohol portion of an ester or the nitrogen atom of the amine portion of the peptide linkage. The transition state can be neither isolated nor detected since it exists for only about 10−13 sec.\nIn molecular terms, these transition states reflect changes in bond lengths and bond angles as well as bond formations and breakages. The energy required to achieve a transition state is denoted as the activation energy which may also be considered as the difference in energy between the energy of the transition state and the energy of the reactants. According to Pauling's explanation, an enzyme preferentially binds the transition state of a reaction, thereby stabilizing it relative to the substrate and products and reducing the activation energy of the reaction, thus increasing the reaction rate. For example, aspartic proteinases are enzymes which are known to catalyze the hyrolysis of peptide linkages within a protein molecule.\nBy extending this explanation, Pauling also predicted that stable analogs of a transition state would bind tightly to an enzyme. In a discussion of substrate distortion as one of several possible sources of rate enhancement by enzymes, it has been suggested that the term “transition state analog” might be used to describe an inhibitor of this kind (1a).\nThe expression of binding energy to compensate for geometrical strain in the substrate is a popular but unlikely mechanism of catalysis. The intermolecular force field cannot overcome the intramolecular force field of the substrate (1).\nCompensation of unfavorable solvation changes by the binding energy is also a possible mechanism of catalysis because of the large solvation energies of the polar groups and ions in water. Lone pairs which may act as general bases or nucleophiles will usually be solvated by hydrogen bonding from water or enzyme. Similarly, other reactive groups will usually be “neutralized” by hydrogen bonding or other mechanisms in the initial state of the enzyme. Usually these groups must be “desolvated” before reaction can occur. This process is energetically expensive and yet an essential part of the normal activation energy which may be compensated by favorable interactions between the non-reacting part of the substrate and the enzyme.\nThe expression of binding energy to compensate for unfavorable electrostatic interactions is another possible mechanism of catalysis. The juxtaposition of like charges could occur in the initial state but be removed in the transition state (2).\nPauling's prediction has become the basis for the now well-established approach to enzyme inhibitor design. The strategy for designing enzyme inhibitors has suggested a strategy for preparing catalytic antibodies whereby antigens are designed based upon mechanistic principles so that antibodies raised in response to such antigens will catalyze a chemical reaction by carrying out the reaction mechanism implicit in the design of the antigen. This strategy has been attempted a number of times.\nFor example, a transition-state analog mimicking an intramolecular 6-member ring cyclization transition state was used to elicit a monoclonal antibody which acted as a stereospecific, enzyme-like catalyst (3). Specifically, the monoclonal antibody so elicited accelerated, by about a factor of 170, the formation of a single enantiomer of a δ-lactone from the corresponding racemic δ-hydroxyester.\nSimilarily, monoaryl phosphonate esters, designated as analogs of the transition state in the hydrolysis of carboxylic esters, were synthesized and used as haptens to elicit specific monoclonal antibodies capable of catalyzing the hydrolysis of carboxylic esters (4). Certain of the antibodies elicited were reportedly found to be catalytic and selective for the hydrolysis of particular aryl esters.\nPhosphonamidates or phosphonate analog-ligands having conformations that substantially correspond to the conformation of a hydrolytic transition state of an amide or ester ligand and which have been used to produce antibodies are described in U.S. Pat. No. 4,659,567 to Tramontano et al. (Tramontano). Antibodies so produced include a paratope that binds to and stabilizes the tetrahedral carbon atom of the amide or ester hydrolysis transition state of the ligand to hydrolyze the ligand at a predetermined site.\nAnalog-ligands which can be used to produce antibody catalysts for the hydrolysis of esters and amides are also described in European Patent Application 0,251,093 of Kollmorgen Corp. (Kollmorgen).\nHowever, none of these analog ligands have been designed in accordance with a rational design approach which maximizes stabilization of the transition state and optimizes atomic relationships within the proteolytic transition state analog, thereby enabling the elicitation of antibodies capable of producing the two dramatic effects of enzyme catalysis; these are molecular recognition and rate acceleration.\nIn enzyme catalysis, groups on both the substrate and the enzyme which are not involved in the chemical mechanism of bond making and breaking, make an important contribution to catalysis. This is illustrated by examining a system where the mechanism is the same for both the enzyme and the non-enzyme catalyzed reaction. For example, the mechanism of action of succinyl-CoA acetoacetate transferase involves nucleophilic attack of the enzyme's glutamate carboxylate on the thioester succinyl CoA, [1] to give an anhydride intermediate (5). The second order rate constant for this reaction is 3×1013 fold greater than the analogous reaction of acetate with the same ester [2]: \nAlthough the nucleophilicities of the carboxylates may be slightly different because of solvation effects and the enzyme may provide some other forms of catalysis these contributions will not be large. The non-reacting part of the enzyme therefore lowers the activation energy by up to 78 KJmol−1. Similarly, changing part of the substrate structure away from the atoms involved in bond making and breaking may also significantly effect catalytic efficiency. The enzyme forms anhydride intermediates from “non-specific” substrates [3]. Even though the chemical reactivities of the two substrates, [1] and [3], are similar, e.g. towards alkaline hydrolysis, the enzyme reaction proceeds up to 3×1012 fold faster with the so called specific substrate [1](5). The non-reacting part of the substrate—the CoA residue—lowers the activation energy by 72 KJmol−1 compared with [3].\nThus, the haptens disclosed in Tramontano do not provide the correct architecture to elicit antibodies that are capable of catalyzing the cleavage of a predetermined peptide sequence in a native protein. These haptens do not provide the correct side-chain groups for production of antibodies that can react with predetermined sites on a protein and cause selective proteolysis in a sequence specific manner. Furthermore, these haptens do not incorporate amino acid side-chain sub-sites on either side of the transition state analog. Without these sub-sites, the haptens cannot provide for the elicitation of catalytic antibodies capable of recognizing a specific amino acid sequence and selectively proteolyzing a peptide linkage within that sequence.\nMoreover, the haptens disclosed in Tramontano do not incorporate leaving groups which facilitate and thereby promote proteolytic catalysis. FIG. 3 (6) depicts a representation of the expected free energy diagram for serine proteinase catalysis. In this mechanism, the enzyme provides a nucleophile (Ser-195). A mechanism involving only a single transition-state requires concerted general base/general acid catalysis by the imidazole (His-57) of the catalytic triad in which the hydrogen from Ser-195 is shared between the serine (Ser 195), imidazole (His-57) and a leaving group, an unlikely situation. Given that a discrete TI exists, there must be two transition states along the acylation pathway (as shown in FIGS. 3 and 4).\nThe stabilization of the carbonyl oxyanion by hydrogen bonds from the protein, the so called “oxyanion hole” in the serine proteinases, is clearly a critical component of the enzyme's catalytic apparatus. The key residue Asn-155 involved in this hydrogen bonding network has been modified by site-directed mutagenesis and results in a decrease of 2–3 orders of magnitude in Kcat, with little change in Km. Decreases in transition state stabilization of 2.2 to 4.7 Kcal mol−1 were observed.\nThe imidazole of His-57 plays two main roles in the acylation step; first, as a base to facilitate attack by the Ser hydroxyl in formation of the tetrahedral species; and second, as an acid in protonating the leaving group in formation of the acyl-enzyme. As shown in FIG. 5, the TI/imidazole system can exist in four ionization states, including the zwitterionic form [2] in which both the imidazole and the carbonyl oxygen are charged. Only the TI species [2] and [3] can directly give transition states leading to the acyl-enzyme. These transition states differ from one another in that the carbonyl oxygen atom is either protonated [3] or in the oxyanion form [2]. Factors involved in determining the actual reaction pathway include the relative pk's of the amine leaving group nitrogen, the serine oxygen and the carbonyl oxygen.\nFor example, an unprotonated amine nitrogen will be a much poorer leaving group than the serine in a tetrahedral transition state (compare pk's of approximately 30 for the former against 15 for the latter). Thus the zwitterionic forms of TI and TS should be less stable, but more efficient at expelling the leaving groups than the corresponding neutral forms. Consequently, the zwitterions are expected to be the catalytically relevant forms (7). Electrostatic interactions within the enzyme's active site will contribute significantly to stabilizing the zwitterionic species.\nIt has been estimated that the pk of the leaving group nitrogen for amide substrates is between 8 and 11 in species such as Tl2 (8). Thus, if the pk of the imidazole remains around 7 in the transition state, then the leaving group has a higher proton affinity than the imidazolium (9). Accordingly, breakdown to the acyl-enzyme would be expected prior to complete formation of the TI. Several lines of evidence, utilizing transition state analog inhibitors, suggest that the true pk of the imidazole in the tetrahedral species is much higher than 7. Although the discussion has centered on the acylation half of the reaction (which is usually the rate-limiting step for amide and peptide substrates) the principles regarding the transition state and tetrahedral intermediate are valid for the deacylation step as well.\nThus, for example, the aromatic compounds described in Tramontano do not lend themselves to provide antibodies that catalyze the hydrolysis of amide bonds in peptides or protein sequences since the pk's of these compounds are significantly more acidic (aniline pka=4.6; p-methylaniline pKa=5.0) and the transition state stabilization afforded by antibodies raised to these immunogens will not provide for hydrolysis of peptide amide bonds since the actual reaction pathways are different.\nMoreover, catalysis occurs through stabilization of the transition state of a reaction. As described below, destabilization of the ground state enzyme/substrate complex relative to the transition state is also necessary for enzymatic catalysis; this prevents accumulation of the enzyme in the enzyme/substrate complex (compare antibody/antigen complex). As noted in Tramontano, it has been suggested that the induced fit mechanism can aid catalysis by providing ground state destabilization relative to the transition state. However, while the induced fit mechanism can increase the value of Km, which reflects the energy of enzyme/substrate complexes, no catalytic advantage is provided for the following reasons.\nThe requirement for destabilization of the enzyme/substrate ground state is illustrated by the reaction catalyzed by two hypothetical enzymes E in FIG. 6 (10,11). Referring to FIG. 6, curve A shows the energetics for the uncatalyzed reaction of substrate S in solution. Curve B shows the reaction when the enzyme stabilizes the ground state and the transition state (ES and ES↑) to the same extent. The enzymatic ground state, ES, is more stable than free S by the amount GB, and the transition state S by the same amount. For the enzymatic reaction depicted by curve B, the rate limiting step is ESEP, and the energetic barrier for this step is ΔGA. This is the same barrier as for the reaction of free S. Thus the enzyme depicted in curve B does not catalyze the reaction, even though the transition state is stablized by this enzyme.\nStill referring to FIG. 6, curve C shows what the enzyme (or antibody) must do; i.e., the transition state complex (ES) must be stabilized as in curve B but the ground state complex (ES) must also be destabilized relative to that in curve B by the amount ΔGB. The energetic barrier for the reaction in curve C, ΔGC, is much less than the barrier for free S (curve A/ΔGA) or the barrier for the enzyme of curve B (also ΔGA). The destabilization of the ground state relative to the transition state that is brought about by the enzyme of curve C is necessary for catalysis. This destabilization corresponds to an increase in Km while the value of Kcat/Km remains constant.\nEnzymatic catalysis will be increased by increasing the energy of the ES complex as long as the transition state energy is not also increased; again, this corresponds to an increase in Km at constant Kcat/Km(12,13). The enhanced efficiency from increasing Km can be explained as an advantage in minimizing the amount of enzyme tied up as ES and thus maximizing the amount of enzyme available for catalysis of the reaction: E+SE+P. The enhancement in the observed rate by increasing Km at constant Kcat/Km is large when the substrate concentration is near or above Km because enzyme, which is complexed with substrate when Km is low, becomes available for catalysis. When the substrate concentration is well below Km a further increase in Km at constant Kcat/Km will yield no advantage because essentially all the enzyme is free and already catalyzing the reaction with the apparent second order rate constant Kcat/Km.\nThe fact that the second order rated constant for reacton of free enzyme and free substrate is Kcat/Km can be derived from the Michaelis-Menten equation:   ν  =                                          K            cat                    ⁡                      [            S            ]                          ⁡                  [          E          ]                    ⁢      tot                      [        S        ]            +              K        m            \nAt low substrate concentration ([S]<<Km), all of the enzyme is free and the Michaelis-Menten equation reduces to ν=(Kcat/Km) [S][E]tot. At high substrate concentration ([S]>>Km), all of the enzyme is in the form ES and the Michaelis-Menten equation reduces to ν=Kcat[E]tot.\nAlthough an increase in Km at constant Kcat/Km is useful for catalysis, the induced fit mechanism increases Km but decreases Kcat/Km, in fact lessening catalysis (13). The decrease is Kcat/Km for the induced fit enzyme relative to the non-induced fit enzyme arises because some of the energy that is used to stabilize the transition state for the non-induced fit must be used to drive the conformational change of the induced fit enzyme. An equivalent explanation is that the decrease is Kcat/Km arises from a lower concentrtion of the active enzyme conformation in the induced fit mechanism (14).\nThe mechanism by which antibodies both destabilize the ground state complex and express intrinsic binding energy for the transition state of the particular reaction to be catalyzed has not been completely understood. An induced fit mechanism for antibody-antigen interaction has been proposed for a monoclonal anti-fluorescein-fluorescein hapten system (15). The binding of the antibody to the hapten was found to display a two-step mechanism; formation of an encounter complex followed by a tightening of binding to the hapten in an “induced fit” mechanism allowing for suitable contact interactions between the antibody combining site and the bound fluorescein ligand. If such a mechanism operates in binding of antibodies to their respective antigenic ligands then there exists an analogy between antibody/antigen interactions and enzyme/substrate interactions.\nIn order for antibodies to function as catalysts, design of the immunogen is crucial for allowing expression of binding energy for the transition state of the reaction. In order for antibodies to destabilize the antibody substrate complex, some form of conformational control of the immunogen is required which will allow for the production of strain and/or distortion of the substrate towards a transition-state trajectory in energetic terms.\nTherefore, the immunogens described in Tramontano and Kollmorgen do not incorporate all these necessary features to provide for antibodies having catalytic function because they have not been rationally designed from knowledge of the mechanistic features of enzyme catalysis. Accordingly, they do not provide suitable templates for generating antibody combining sites endowed with catalytic properties."}
{"text": "Field of the Disclosure\nThis application relates, in general, to wake control systems for a watercraft, and some embodiments relate to surf wake systems for modifying a wake produced by a watercraft travelling through water.\nDescription of the Related Art\nWake surfing has become increasingly popular in recent years because, unlike an ocean wave, a wake produced by a watercraft is on-demand not to mention continuous and endless as long as the watercraft is moving forward. As a watercraft travels through water, the watercraft displaces water and thus generates waves including bow wave and diverging stern waves on both sides of the watercraft. Due to pressure differences, these waves generally converge in the hollow formed behind the traveling watercraft and/or interfere with each other to form a wake behind the watercraft. Such a wake, however, is generally small, choppy or too close to the watercraft to be suitable and safe for water sports, and particularly not suitable for wake boarding or surfing.\nTo facilitate surfing, a wake can be formed away from the stern of the watercraft, for example, about ten feet away, and with a waist-height peak, for example, about three feet or higher. Those of skill in the art will understand that a wake for wake surfing can be formed at various different distances behind the watercraft, and the wake can have various different heights. Generally hundreds, and sometimes thousands, of pounds of additional weight or ballast to a rear corner of the watercraft to make the watercraft tilt to one side, displaces more water, and hence generates a larger wake on that side. Such additional weight may be in the form of removable ballast bags, installed ballast tanks or bladders, or passengers positioned to one side of the watercraft, which is primarily used to tip the watercraft to that side. Using such additional weight to produce larger wakes, however, poses several disadvantages. For example, such additional weight may take up significant space and capacity that may otherwise reduce the passenger capacity of the watercraft. Also, such additional weight may unbalance the watercraft creating difficulties in control. Moreover, the additional weight generally must be moved from one side of the water craft to the other in order to generate a wake on the other side of the water craft. Shifting such additional weight may require significant time and effort. For example, filling and emptying ballast tanks to switch from one side to the other may require 20 minutes or more.\nIn light of the foregoing, it would therefore be useful to provide surf wake system that overcomes the above and other disadvantages."}
{"text": "In well drilling, particularly for oil or gas, casing is cemented in the well to seal off the formation. For pressure control and the protection of fresh water formations, the well will have multiple strings of casing, each string of casing being smaller in diameter than the one above. In one technique, each string of casing extends completely to the surface where it is supported at a wellhead by a casing hanger. In another technique, liners are employed. A liner is a string of casing that will overlap the next upward string of casing, but will not extend completely to the surface. This avoids additional costs of several strings of concentric casing in the shallower portions of the well.\nThe upper end of the liner has to be remotely attached by a liner hanger to the casing, normally near the lower end of the casing. When installing a liner, normally a special section of pipe, called a mandrel, will be attached to the upper end of the liner. The mandrel will be machined to accept the liner hanger assembly. Because the mandrel must have a performance rating equal to the pipe of the liner, normally the customer will supply for machining a section of the actual liner to be used. A running tool will engage the mandrel, and the entire assembly will be lowered into the well. After the liner reaches bottom, the operator will actuate the liner hanger to secure the mandrel to the casing.\nThe operator will cement the liner into the well, in some cases before actuating the liner hanger, on other cases after. Usually a cement bushing will be employed with the running tool to locate inside the mandrel. This cement bushing engages the mandrel to support the running tool and drill pipe against an upward force that occurs while pumping the cement down the drill pipe. After the cement has set and the liner hanger has been set, the operator withdraws the running tool and cement bushing. The mandrel and liner hanger assembly will remain in the well.\nThere are many different types of liner hangers and different techniques employed for setting liner hangers. In one technique, hydraulic force is used to set the liner hanger. The hydraulic force may be supplied by pumps on the rig pumping down the drill pipe. The flow must be diverted through a fluid port in the running tool and the mandrel to apply the fluid pressure to the liner hanger on the exterior of the mandrel. The flow is diverted by dropping a ball or dart from the surface, which opens and closes various passages once seated. A considerable amount of time is required for the ball to reach the running tool, often several hours. During this time period, the operator is normally unable to circulate fluid through the liner. If the liner has not already been cemented in place, keeping the liner stationary in an open hole without the ability to circulate can be dangerous.\nIn another hydraulic type, hydraulic force is supplied by a hydraulic chamber exposed to wellbore fluids. The operator must raise and lower the drill pipe to provide the increased pressure. This system may prevent the operator from raising and lowering the drill pipe prior to setting for other purposes.\nOther liner hangers are operated mechanically, which involves rotating and moving the drill pipe axially to actuate the liner hanger. For example, J-pins and J-slots may be employed which, when rotated, move the liner hanger to different positions. In highly deviated wells, it is difficult to rotate the drill pipe. Other mechanical running tools do not have the ability for the operator to stroke the liner up and down after the liner reaches bottom without causing setting of the hanger. As a result, the operator may be unable to reciprocate the liner during cementing, which often is a good procedure to follow.\nThe requirement of specially prepared mandrels is a disadvantage. The machining is done at a shop, not in the field. Consequently, delays may occur, and it is costly to ship the liner section to a shop and back to a wellsite. Moreover, in some costly corrosion-resisting casing strings, machining may damage the properties of the mandrel. Inspections are necessary to make sure that the machining operations have not harmed the metal of the pipe.\nAdditionally, in very long liners it may be desired to spread the load support by the liner hanger to the casing at various points along the overlapped region. Liner hangers used, however, support the liner only at the top of the liner. Other liner hanger systems and running tools may not have the ability to automatically fill the drill pipe while the liner is being run. The customer may prefer to have the drill pipe filled while being run. If so, then the customer must manually fill it from the drill rig, which takes additionally time."}
{"text": "This invention relates to bioadhesive compositions, particularly electrically conductive hydrogel compositions having bioadhesive properties. The invention further relates to biomedical skin electrodes incorporating such hydrogel bioadhesive compositions that are electrically conductive.\nBiomedical skin electrodes are widely used in a variety of situations, whenever for example it is required to establish an electrical connection between the surface of the body of the patient and external medical equipment for transmission of electrical signals.\nModem medicine uses many medical procedures where electrical signals or currents are received from or delivered to a patient\"\"s body. The interface between medical equipment used in these procedures and the skin of the patient is usually some sort of biomedical electrode. Such electrodes typically include a conductor which must be connected electrically to the equipment, and a conductive medium adhered to or otherwise contacting skin of the patient, and they are of varying types with a wide variety of design configurations which will generally depend on their intended use and whether for example they are to be used as transmission electrodes or sensing i.e. monitoring electrodes.\nAmong the therapeutic procedures using biomedical electrodes are transcutaneous electric nerve stimulation (TENS) devices used for pain management; neuromuscular stimulation (NMS) used for treating conditions such as scoliosis; defibrillation electrodes to dispense electrical energy to a chest cavity of a mammalian patient to defibrillate heart beats of the patient; and dispersive electrodes to receive electrical energy dispensed into an incision made during electrosurgery.\nAmong diagnostic procedures using biomedical electrodes are monitors of electrical output from body functions, such as electrocardiograms (ECG) for monitoring heart activity and for diagnosing heart abnormalities.\nFor each diagnostic, therapeutic, or electrosurgical procedure, at least one biomedical electrode having an ionically conductive medium containing an electrolyte is adhered to or is otherwise contacted with mammalian skin at a location of interest and is also electrically connected to electrical diagnostic, therapeutic, or electrosurgical equipment. A critical component of the biomedical electrode is the conductive medium which serves as the interface between the mammalian skin and the diagnostic, therapeutic, or electrosurgical equipment, and which is usually an ionically conductive medium.\nBiomedical electrodes are used among other purposes to monitor and diagnose a patient\"\"s cardiovascular activity. Diagnostic electrodes are used to monitor the patient immediately and are only applied to the patient for about five to ten minutes. Monitoring electrodes, however, are used on patients in intensive care for up to three days continuously. In contrast, Holter electrodes are used to monitor a patient during strenuous and daily activities.\nAlthough all of the biomedical electrodes just referred to are used to record cardiovascular activity, each electrode requires specific features or characteristics to be successful. Thus, the diagnostic electrode does not have to remain adhered to a patient for extensive periods but it does have to adhere to hairy, oily, dry and wet skin effectively for the five to ten minutes of use. The monitoring electrode has to adhere for a longer period of time although the patient is often immobile during the monitoring period. The Holter electrodes is susceptible to disruption from adhesion due to physical motion, perspiration, water, etc., and therefore requires the best adhesion and at the same time comfort and electrical performance.\nIn the biomedical electrodes known in the prior art the ionically conductive medium which serves as an interface, between the skin of a mammalian patient and the electrical instrumentation, ranges from conductive gels and creams to conductive pressure sensitive adhesives. However, while the conductive media can be in the form of pressure sensitive conductive adhesives, for monitoring or Holter biomedical electrode use such conductive adhesives are not generally adequate on their own to maintain adhesion to mammalian skin and additional hypoallergenic and hydrophobic pressure sensitive adhesives may be employed around the conductive medium to provide the required mammalian skin adhesion. U.S. Pat. No. 5,012,810 (Strand et al.) and U.S. Pat. Nos. 4,527,087, 4,539,996, 4,554,924 and 4,848,353 (all Engel) are examples of documents that disclose biomedical electrodes which have a hydrophobic pressure sensitive adhesive surrounding the conductive medium.\nIn general, a desirable skin electrode is one which maintains good electrical contact with the skin and is free of localised current hot spots, i.e. exhibits uniform conductivity. For example, it has been found that a prior art electrode utilising karaya gum tends to creep in use and flatten out, exposing skin to possible direct contact with the current distribution member or lead wire. A desirable skin electrode should also usually have a low electrical impedance.\nIt is an object of this invention to provide hydrogel adhesives possessing controlled and predictable adhesive properties which may be readily varied to suit different uses and, in the case of medical electrodes or similar devices, different configurations or applications. It is also an object of the invention to provide such hydrogel adhesives which in addition may possess superior electrical characteristics as compared to those commonly associated with bioadhesive hydrogels."}
{"text": "The present invention relates to a write control device, and more specifically, to a technology for rapidly performing a write operation of a nonvolatile memory device configured to sense data based on resistance change.\nGenerally, memory devices may be classified as volatile memory devices or nonvolatile memory devices. A nonvolatile memory device includes a nonvolatile memory cell capable of preserving stored data even when a power source is off. A nonvolatile memory device may be implemented, for example, as a flash random access memory (RAM) device, a phase change RAM (PCRAM) device, or the like.\nPCRAM devices include memory cells that are implemented using a phase change material, for example, germanium antimony tellurium (GST), and are configured to store data in the memory cells by applying heat to the GST so that the GST changes into a crystalline phase or an amorphous phase.\nA nonvolatile memory device, such as a magnetic memory device, a phase change memory (PCM) device, or the like, has a data processing speed similar to that of a volatile RAM device. However, unlike a volatile RAM device, a nonvolatile memory device preserves data even when a power source is off.\nFIGS. 1a and 1b illustrates a conventional phase change resistance device 4.\nReferring to FIGS. 1a and 1b, a conventional phase change resistance device 4 includes an upper electrode 1, a lower electrode 3, and a phase change material 2 interposed between the upper electrode 1 and the lower electrode 3. When a voltage is applied to the upper electrode 1 and the lower electrode 3, a current flows into the phase change material 2, thus inducing a high temperature in the phase change material 2. As a result, the electrical conductive state of the phase change material 2 changes depending on resistance variation due to the high temperature.\nFIGS. 2a and 2b illustrate a phase change principle of the conventional phase change resistance device 4.\nReferring to FIG. 2a, if a low current smaller than a critical value flows into the phase change resistance device 4, the phase change material 2 is crystallized. When the phase change material 2 changes into a crystalline phase, it becomes a low resistance material. As a result, a current can flow between the upper electrode 1 and the lower electrode 3.\nOn the other hand, referring to FIG. 2b, if a high current greater than the critical value flows into the phase change resistance device 4, the phase change material 2 has a temperature higher than a quenching point. When the phase change material 2 changes into an amorphous phase, i.e., a non-crystalline phase, it becomes a high resistance material. As a result, a current cannot easily flow between the upper electrode 1 and the lower electrode 3.\nThe phase change resistance device 4 can store data corresponding to two resistance phases. That is, in one case, if a low resistance phase in the phase change resistance device corresponds to a data “1,” and a high resistance phase corresponds to a data “0,” then the phase change resistance device 4 may store two logic states of data.\nThis data can be stored in the phase change resistance device 4 as nonvolatile data because the status of the phase change material 2 does not change even when a power source is off.\nFIG. 3 illustrates a write operation of a conventional phase change resistance cell.\nReferring to FIG. 3, heat is generated if a current flows between the upper electrode 1 and the lower electrode 3 of the phase change resistance device 4 for a given time.\nIf a low current, smaller than the critical value, flows for a given time, the phase change material 2 changes into a crystalline phase. As a result, the phase change resistance device 4 becomes a low resistance element having a set phase.\nOn the other hand, if a high current, greater than the critical value, flows for a given time, the phase change material 2 changes into an amorphous phase. As a result, the phase change resistance device 4 becomes a high resistance element having a reset phase.\nAccordingly, a low voltage is applied to the phase change resistance device 4 for a long period of time in order to write the set phase in the write operation.\nOn the other hand, a high voltage is applied to the phase change resistance device 4 for a short period of time in order to write the reset phase in the write operation.\nTo change the phase change resistance cell into the set phase, it is important to control a quenching slope of a set write current, which is required for crystallizing the phase change resistance cell, by gradually reducing an amount of the set write current. This way of gradually reducing the set write current is called “quenching.”\nHowever, for example, if a reference current, which is used to generate the set write current and is received from the outside or generated inside the chip, changes by some factors or if a clock having a wrong value is generated by mismatched circuits, the set write current may be generated with an undesired quenching slope.\nMoreover, without accurately checking the quenching slope, it is impossible to precisely control the phase change resistance cell in program and verify (PNV) operations and a multi-level cell MLC where multi-leveled resistance distribution of a phase change material, e.g., germanium antimony tellurium (GST), is formed.\nMeanwhile, a conventional nonvolatile memory device controls a driving voltage supplied to a memory cell array in a program operation using a write driving unit.\nThe write driving unit is configured to stably supply a write current even though the write current has a small current value. However, it takes a much longer time to supply cells with the write current when using a small current value than when using a large current value.\nThat is, when the write driving unit supplies a small write current to the cells, it takes much longer to supply the current because of parasitic components on a write path.\nFIG. 4 illustrates a write time that is delayed by a large parasitic component and a small write current. Here, IT represents a driving current flowing in a driving transistor of the write driving unit, and IM represents a cell current flowing in a memory cell.\nReferring to FIG. 4, the driving current IT flowing in the driving transistor reaches a stabilized current level after a settling time tsettling passes. In addition, there is a time delay due to parasitic components on a write path. A time delay tDELAY occurs in a period of time where the driving current IT of the driving transistor is transmitted to the memory cell.\nA small write current is generally determined by characteristics of cells and the size of a cell array. As the size of the cell array becomes larger, the delay time occurring due to the parasitic components on the write path also become greater.\nIn this way, the write time becomes longer by the small write current and the large size of the cell array. A conventional write driving unit has not introduced improvements to such a write time delay phenomenon."}
{"text": "1. Field of the Invention\nThis invention pertains to a resolver to digital converter and more particularly to a resolver to digital converter which uses digital signals to provide accurate control of the excitation voltages applied to the resolver windings so that the accuracy of the digital output representation is improved.\n2. Description of the Prior Art\nIn many types of equipment it is desirable to reliably and accurately measure the position of a rotatable shaft and to convert the resulting measurement into a digital signal which is representative of the shaft position. One such type of equipment where it is desirable to digitally measure angular position is a weapons delivery system where the elevation and azimuth pointing angles must be accurately measured. Some of the prior art converters make use of the conventional resolver or synchro transmitter having a rotor winding and a plurality of stationary windings. A voltage source and a phase shifter supply multiphase AC voltages to the stator windings to induce a magnetic flux which rotates at the frequency of the applied voltage and induces an AC voltage in the rotor winding. The AC voltage in the rotor winding has a phase angle which varies with the angular position of the rotor. Thus, the angular position of the rotor can be determined by measuring the difference in phase angle between the voltage induced in the rotor winding and a fixed reference voltage. Such difference can be measured by determining the time which elapsed between the zero passage of the reference voltage and the zero passage of voltage from the rotor winding, and this time can be expressed as a digital representation of the rotor angle. A serious disadvantage of the foregoing type of equipment is that a very high degree of accuracy must be maintained between the phase of the voltages applied to the various stator windings. Aging of the components used in the circuitry may cause a change in phase between the voltages applied to the various stator windings and therefore cause inaccuracies in the measured angle between the reference voltage and the voltage produced by the rotor winding. These inaccuracies lead to an inaccurate output reading of the angular position of the rotor.\nOne type of prior art apparatus as shown in prior U.S. Pat. No. 3,634,838, uses a high-frequency signal source to develop signals for the stator windings of a resolver or goniometer having three stator windings and a rotor winding. The signal source develops high-frequency pulses which are reduced to low-frequency pulses by a frequency divider. The output of the frequency divider is coupled to a filter which converts the pulses into a sinusoidal voltage that is coupled to one of the stator windings. The sinusoidal voltage from the filter is applied to a first stator winding and to a 90.degree. phase shifter. The phase shifter develops a second sinusoidal voltage shifted 90.degree. from the first voltage, the second sinusoidal voltage being applied to a second stator winding. A counter which is coupled to the signal source is used to count the high-frequency periods. In addition, a phase detector is provided for comparing the phase of an output voltage from the counter with that of a reference voltage from the goniometer. The phase of the pulses from the frequency divider is controlled relative to the phase of the voltage from the counter by varying the number of pulses supplied to the frequency divider. This prior art apparatus has no provision for correcting any variation in the phase angle between the voltages applied to the first and second stator windings of the goniometer."}
{"text": "In a conventional electromagnetic relay described in JP-A-2008-226547, two fixed contact supports having two fixed contacts respectively are positioned, and two movable contacts are contacted by and separated from the two fixed contacts by moving one movable body, to which the two movable contacts are fixed, so that an electrical circuit is opened and closed. More specifically, the electromagnetic relay includes a movable member driven by electromagnetic force of a coil, a pressure spring for urging the movable body such that the movable contacts contact the fixed contacts, a return spring for urging the movable body via the movable member such that the movable contacts are separated from the fixed contacts, and the like. Further, a contact surface between the movable contacts and the fixed contacts is a spherical surface.\nWhen the coil is energized, the movable member is driven toward a side away from the movable body by the electromagnetic force, and the movable body is urged by the pressure spring to be moved. As a result, the movable contacts contact the fixed contacts, and the movable member is separated from the movable body.\nIn a conventional electromagnetic relay described in JP-A-62-51126, two fixed contact supports having two fixed contacts respectively are positioned, and two movable contacts are contacted by and separated from the two fixed contacts by moving one movable body, to which the two movable contacts are fixed, so that an electrical circuit is opened and closed. More specifically, the movable body is integrated with a movable member that is driven by electromagnetic force of a coil such that the movable body can be relatively moved, the movable body is held at a predetermined position of the movable member by a pressure spring, and the movable member and the movable body are urged by a return spring such that the movable contacts are separated from the fixed contacts. Further, a contact surface between the movable contacts and the fixed contacts is a spherical surface.\nWhen the coil is energized, the movable member and the movable body are driven by the electromagnetic force and the movable contacts contact the fixed contacts. At this time, the pressure spring is bent by a stroke of the movement of the movable member after the movable contacts contact the fixed contacts, thereby the movable member and the movable body are relatively moved.\nHowever, in the electromagnetic relay described in JP-A-2008-226547, because the contact surface is a spherical surface, the movable contacts contact the fixed contacts by point contact. When the coil is energized, the movable contacts contact the fixed contacts by two-points contact, that is, at a first contact portion between one fixed contact and one movable contact, and a second contact portion between the other fixed contact and the other movable contact. In this manner, because the movable contacts contact the fixed contacts by the point contact, i.e., the two-points contact, the movable body may vibrate around a line passing through the first and second contact portions when the movable contacts collide with the fixed contacts. In particular, in the electromagnetic relay described in JP-A-2008-226547, because the movable member is separated from the movable body when the movable contacts contact the fixed contacts, the vibration suppression effect by the movable member cannot be obtained, and thereby it becomes difficult to suppress the vibration of the movable body. Thus, the vibration of the movable body may be resonated in a casing to generate abnormal noise.\nIn the electromagnetic relay described in JP-A-62-51126, because the movable body is integrated with the movable member, it is easy to suppress vibration of the movable body when the movable contacts collide with the fixed contacts. However, a position of a contact portion is changed in accordance with the vibration of the movable body until the vibration of the movable body is suppressed. Thus, a resistance value between the contacts is changed, and thereby it becomes easy for the contact portion to generate heat and the wear-and-tear of the contacts may occur."}
{"text": "The invention relates to an improved side-guiding mechanism for a folding machine. In particular, the invention relates to an improved side guiding mechanism for a folding machine which makes outserts.\nAn outsert is a sheet which is folded a number of times with the outermost fold being glued tight to an inner fold in order to make a tightly folded packet. FIG. 1 illustrates a typical outsert 12 wherein a last wrap-around panel 14 is adhered by a glue or adhesive spot 15 to an adjacent inner fold 16 to complete outsert 12. Outserts are often very long, for instance, 9 to 50 inches in length, and are often used by the pharmaceutical industry to provide information concerning the drug to which the outsert is attached. Such information might include instructions for use or warnings about possible side effects of the drug. Generally, they are attached to the bottle or box in which the drug is packaged.\nU.S. Pat. No. 4,812,195 (incorporated herein by reference) describes an apparatus for automatically forming sheets from a web of material. The apparatus described therein includes an outsert attachment which is attached to the rest of the folding apparatus. The outsert attachment folds and glues the sheets in such a manner as to form outserts similar to the one shown in FIG. 1. FIG. 2, which is adapted from U.S. Pat. No. 4,812,195, illustrates this prior art apparatus for forming outserts of a predetermined length from a web 20 by accumulating long lengths, for example, up to 50 inches of the web 20, in an accumulator 28. A photocell 29 detects marks on the web 20 to operate a severing means 31 at a sheet station 26 to form sheets 11 each of which is the same size. A sheet conveyor 35 extends from the sheet station 26 to a folding station 33. At the folding station 33, folding plates 37 and rollers.(not shown) fold the sheets 11 into the form shown on conveyor 38. As can be seen, a sheet 11 at this point has a folded portion 39 and an unfolded wraparound panel 14, with the folded portion 39 leading the wrap-around panel 14 in the direction of travel shown by the arrow. The conveyor 38 conveys the sheet in this fashion into a wraparound and gluing station 40 of the outsert attachment.\nThe precision of the last fold, and the gluing of the wraparound fold 14, is accomplished by precisely positioning the folded sheet at the wraparound and gluing station 40. More specifically, the leading edge of the folded portion 39 is stopped in its forward travel by stops 41 and, in a fraction of a second of its stopping, a side jogging mechanism 42 is activated to push the folded sheet against a side stop 43 to that the sheet is positioned precisely. As soon as this occurs, an adhesive applicator 45 is actuated to apply one or more adhesive spots 15 to the folded section 39. Immediately following this, a folding means 48 which includes an overhead folding knife 49 is operated to fold the wraparound fold 14 against the adhesive spots 15 on the inner fold 16 thereby adhering the folds together and completing the formation of the outsert.\nTypically, the sheet is printed with several outserts side-by-side which must then be separated. To accomplish this, the overhead folding knife 49 which engages the folded section 39 of the folded sheets pushes the folded section 39 into the nip of underlying folding rollers 50. The folding rollers 50 force the wraparound fold 14 tightly against the glue spots 15 and against the inner fold 16 to complete the outsert. The folding rollers 50 are also provided with circular knives 53 mounted on a pair of slitting rollers. 52 which slit the sheet into a multiplicity of side-by-side outserts. The outserts are then re-oriented and discharged automatically in an edge position into a container 22.\nFIG. 3, which is also adapted from U.S. Pat. No. 4,812,195, illustrates in detail the side guiding mechanism used in the apparatus described therein. The illustrated side jogger 42 comprises a solenoid 150, slidably mounted on a support slide bracket 151 mounted on the inside of a stationary side frame member 156. The support bracket 151 has a slide channel 152 into which is mounted slide 153. Suitable lock fasteners 154 lock the slide 153 in position on the slide bracket 151 to position the solenoid plunger 158. Thus, the solenoid 150 can be positioned to accommodate various widths of sheets to be side jogged. The solenoid 150 includes a plunger 158 with a sheet pusher 159 mounted at its end. The sheet pusher 159 abuts against the side edge 39a of the folded sheet to push the same laterally, thereby bringing an opposite edge 39b of the folded sheet against the side stop 43 while the sheet slides along and against the forward stops 41. The solenoid 150 includes a return spring 161 extending between a head 162 on the plunger and the solenoid body 163. A stop nut 164 on the other threaded end 165 of the solenoid limits the return travel by the spring 161.\nThe operation of the side jogger is as follows. The plunger/pusher 158/159, when activated by electronic sensing means, pushes the mostly folded paper against the side stop 43. This aligns (i.e., side-guides) the folded sheet for positive and consistent gluing and slitting. As the width of the mostly folded paper varies between different jobs in a typical range of 7 to 9 inches, and as the width of the paper may vary by about 1/16 of an inch during the course of a single job, a simple and safe method of adjusting both the plunger/pusher 158/159 and the side stop 43 is essential for the efficient and safe operation of the outsert attachment.\nIn the apparatus of U.S. Pat. No. 4,812,195, the plunger 158 extends through the body of the solenoid 150, and when the solenoid is electrically activated, the pusher 159 mounted on the end of plunger 158 travels approximately 1/4 inch to 1/2 inch. The pusher/plunger is then returned to its original position by means of spring 161. The whole assembly is bolted directly onto the side frame member 156 of the outsert attachment. While the apparatus of U.S. Pat. No. 4,812,195 discloses that adjustment is possible by sliding the entire side jogging mechanism in slide bracket 151, the commercial embodiment of this apparatus is sold without this slide mechanism. In commercially available machines, the adjustment is maintained by the use of the locknut 154 and associated washers probably because the slide mechanism is unable to maintain the desired position of the pusher/plunger during the course of the run. The locknut and washers used in the commercial embodiment permit only a limited amount of adjustment for different jobs. The locknut and washers serve a dual purpose as they also regulate the amount of the plunger's return travel to the side frame member. Because the travel distance is severely limited, the operator must maintain an inventory of different size plungers and pushers to accommodate the various widths of outserts for different jobs.\nWhen setting up the outsert attachment for actual operation, the proper combination of plungers and pushers has to be selected and then adjusted. At the same time, the position of the side stop also has to be adjusted. The stop is bolted directly onto the folding machine which is mounted inside the side frame members of the outsert attachment. To adjust the stopper, a set screw (not shown) has to be loosened and then the stopper can be slid to the desired position. The set screw is then tightened to hold it in place. As the stopper slides freely, finding a proper position can be very difficult. This is especially true since the stopper often slips out of position and requires further adjustment when the set screw is tightened. The entire process of finding the proper plunger/pusher combination and adjustment, and the proper stopper adjustment can take as long as an hour in an initial machine setup. Furthermore, during the course of a run, constant adjustment is required as the width of the mostly folded sheet may vary by as much as 1/16 of an inch. These adjustments add hours to the time needed to complete a run.\nAdditionally, because the above-described adjustment process is so cumbersome, it is often carried out in a dangerous manner. As it is quite difficult to make the proper adjustments while the machine is idle, many operators are tempted to make adjustments \"on the fly,\" i.e., while the machine is running. When the adjustments are made on the fly, the operator can see very well if the outserts are being side-guided properly. However, the problem in making adjustments on the fly is that there are many moving metal parts in the immediate proximity of the pusher/plunger and side stop, including the circular slitting knives. It is exceedingly dangerous to make any adjustments at all while the slitting knives are in operation. Many accidents have occurred in this manner.\nAnother disadvantage of mounting the side stop directly onto the folding apparatus is that from time to time, the outsert attachment has to be removed from the rest of the machine for repair during a job. In such cases, the side stop also has to be removed and then re-installed when the outsert attachment is replaced. This entails the re-adjustment once again of the side stop with all the inconveniences associated therewith.\nAccordingly, it is an object of the present invention to provide an outsert attachment which has a more efficient and a safer mechanism for adjusting the plunger/pusher and side stop combination."}
{"text": "Silicon is a common semiconductive material used in the fabrication of integrated circuits. Silicon can occur in crystalline and amorphous forms, and when crystalline can be monocrystalline or polycrystalline. In some instances, silicon is combined with germanium, essentially forming a silicon germanium alloy. Such materials can be doped with conductivity enhancing impurities (i.e., boron and/or phosphorus) to modify the conducting characteristics of the silicon-comprising material.\nMonocrystalline silicon can be provided in bulk substrate form or otherwise grown or deposited epitaxially from an exposed monocrystalline material. Epitaxy generally involves the growth or deposition of a single or monocrystalline layer of material such that the epitaxial layer has a crystal orientation which is common to that of the material from which it is grown. One factor that determines the quality of the epitaxial silicon-comprising layer relates to the presence and quantity of crystallographic defects. Such are non-uniformities in the crystal structure of the epitaxial layer. Many of these defects are caused by defects appearing at the surface of the substrate which propagate into the layer during growth. Examples include sidewall dislocations and stacking faults. Dislocations and stacking faults can be electrically active more so than the surrounding material within the epitaxial layer due to the presence of dangling bonds. Such can lead to unnecessary recombination generation currents, lower breakdown voltages, higher current leakage and larger junction ideality factors.\nOne place where epitaxial silicon has been utilized is as one or more of the components in a field effect transistor. Transistor structures comprise a channel region received between a pair of source/drain regions, and a gate configured to electrically connect the source/drain regions to one another through the channel region. The transistor constructions utilized in semiconductor constructions are supported by a semiconductor substrate. The semiconductor substrate will have a primary surface which can be considered to define a horizontal direction. Transistor devices can be divided into two broad categories based upon the orientations of the channel regions relative to the primary surface of the semiconductor substrate. Specifically, transistor structures which have channel regions that are primarily parallel to the primary surface of the substrate are referred to as planar or horizontal transistor structures, and those having channel regions which are generally perpendicular to the primary surface of the substrate are referred to as vertical transistor structures. Since current flow between the source and drain regions of a transistor device occurs through the channel region, planar transistor devices can be distinguished from vertical transistor devices based upon the direction of current flow as well as on the general orientation of the channel region. Specifically, vertical transistor devices are devices in which the current flow between the source and drain regions of the devices is primarily substantially orthogonal to a primary surface of a semiconductor substrate, and planar or horizontal transistor devices are devices in which the current flow between source and drain regions is primarily parallel to the primary surface of the semiconductor substrate.\nEpitaxial silicon-comprising materials have been proposed for use in channel regions of vertical transistors. Further, one or both of the source/drain areas of a vertical transistor might also comprise epitaxially grown silicon or an epitaxially grown silicon germanium alloy. Requirements for epitaxial materials within a vertically oriented channel region are typically more stringent than for the use of such material in source/drain regions of horizontally oriented field effect transistors. Further, fabrication of vertical field effect transistors typically utilizes masks of oxide, nitride or other materials for self-aligned patterning of the epitaxial silicon-comprising material during its formation. The interface of the epi with these materials can be a defect source. Further, the selective epitaxial growth of silicon for vertical transistors typically utilizes lower deposition temperatures as compared to blanket epitaxial silicon depositions. Unfortunately, the use of lower temperatures reduces surface mobility and can also result in increased defects over that of higher temperature processing. Also and regardless, thermal stress can be generated during cool-down of the substrate from the temperature at which the epitaxial silicon-comprising material was grown. This can result in crystallographic defects being generated after growth.\nFurther, where the epitaxial silicon-comprising material includes germanium, such has an increased tendency for defect formation on the surface during deposition due to mismatched lattice constants of silicon and germanium. These defects propagate and either terminate with other defects or at the surface. Regardless, after deposition, crystallographic defects are extremely difficult to remove or heal within the bulk epitaxially grown material or at interfaces of such material with other materials.\nWhile the invention was motivated in addressing the above identified issues, it is in no way so limited. The invention is only limited by the accompanying claims as literally worded, without interpretative or other limiting reference to the specification, and in accordance with the doctrine of equivalents."}
{"text": "Polymeric aromatic carbamic acid esters (polyurethanes) such as diphenylmethane dicarbamates and the related higher homologs polymethylene polyphenyl carbamates have become increasingly important products particularly, for use in the preparation of the commercially valuable diphenylmethane diisocyanates and mixtures of diisocyanates and the polyisocyanates by the decomposition of such polymeric aromatic carbamic acid esters in a suitable solvent as shown in Rosenthal et al, U.S. Pat. Nos. 3,962,302, June 8, 1976 and 3,919,279, Nov. 11, 1975.\nAt the present time there is no known successful commercial method for the direct preparation of polymeric aromatic esters of carbamic acid. The corresponding diphenylmethane diisocyanates and polyisocyanates, available commercially, are largely produced by the phosgenation of mixtures of diamines and polyamines obtained by the condensation of aniline and formaldehyde with catalytic quantities of a mineral acid as for example, disclosed in the Pistor et al, U.S. Pat. No. 4,014,914.\nA proposed prior art process for the preparation of polymeric aromatic carbamic acid esters (polyurethanes) is disclosed in Klauke et al, U.S. Pat. No. 2,946,768 and involves the condensation of aryl carbamic acid esters with carbonyl compounds in a dilute aqueous acid condensation medium. However, in such process the carbonyl compound such as formaldehyde tends to react at the nitrogen of the carbamate to produce, along with desired polyurethanes, varying amounts, i.e., generally between 15 percent and 50 percent by weight, of undesirable (alkoxycarbonyl)phenylaminomethylphenyl compounds which includes the various dimers, trimers, tetramers, etc. of such compounds (also referred to herein as \"N-benzyl\" compounds). Attempts to prepare mono or diisocyanates and polyisocyanates or to otherwise use the mixture containing the undesired N-benzyl compounds, which cannot be converted to an isocyanate by pyrolysis, and polyurethanes presents many problems since there is no known method for separating the polyurethanes from the N-benzyl impurities."}
{"text": "A multitude of security devices exists within the prior art which teach the use of supplementary devices in conjunction with a typical latching arrangement for preventing undesirable access to one's premises. Examples of such devices are supplementary chains, dead bolts and latches.\nU.S. Pat. No. 4,484,463 issued on Nov. 27th, 1984 by Hennessy describes a door lock guard for protecting the cylinder of a door lock upon the key cylinder side. Such door lock is permanently mounted as best illustrated in FIG. 3 by fastening elements 46. The shape of portion 18 of plate 12 formed from the appropriate materials ensures that the cylinder 34 cannot be accessed by a tool and subsequently rotated.\nU.S. Pat. No. 4,074,552 describes a kit which adapts to a lock preventing tampering of the lock thus undesirable access to one's premises. Such a combination is accessible only by a special long key which will turn the cylinder of the lock.\nU.S. Pat. No. 3,556,571 by Laub issued Jan. 19, 1971 describes a device for attachment to a door knob which in essence child proofs the door knob and prevents rotation of the lock bolt of a door until such time as an adult removes the attachment. The instant invention provides a number of alternative embodiments of a latch blocking element which is installed upon the interior of a door proximate the operating means of a door latch and thus prevents the rotation of the operating means in use until such time as the blocking device is repositioned. The blocking device thus prevents the rotation of the for example a key cylinder in a household door. Should the key cylinder be tampered with on the exterior of the door by an intruder, the rotation thereof would be prevented on the interior by blocking the rotation of the operating means connected to the key cylinder. Nowhere within the prior art is such a blocking means found which may be used by for example the occupant of a home as a night latch for securing his or her home.\nIt is therefore an object of this invention to provide a latch blocking mechanism which blocks the operation of operating means of a latching means thus preventing the rotation of the latching means from a locked to an unlocked position.\nIt is a further object of this invention to provide a blocking means embodied in alternative embodiments which cooperate with those latching means most broadly utilized in household latches. It is a further object of the invention to provide a latch blocking means which is easy to install, move, remove, or reposition.\nFurther and other objects of the invention will become apparent to a man skilled in the art when reviewing the following summary of the invention and the more detailed description of the preferred embodiments illustrated herein."}
{"text": "This invention relates to methods of evaluating the condition of oilfield tubular goods. More specifically, this invention relates to methods and apparatus for determining by measurement the extent of internal wear of subsurface pipe or casing.\nA variety of problems may result from internal casing wear. Drill pipe collars may rub the casing, and possibly lead to a blowout in the casing if extremely worn. In a production well, badly worn casing may permit the undesirable \"thiefing\" of the flow to unwanted zones and thereby reduce surface production. In an injection well, worn casing may permit the injected fluid to flow to undesired formations.\nInternal wear, which is the removal of metal from the inside casing wall occurs in three principal patterns: loss of metal from only part of the inside casing wall while the rest is unworn; uniform loss of metal around the entire inside casing wall circumference; and non-uniform loss of metal around the entire inside casing wall. In addition, there are two patterns associated with unworn casing--round internally with all of the original metal in place, and elliptical with all of the metal in place.\nA number of surveys are available for assessing internal wear in downhole casing. One such survey provides a depth recording which responds to the circumferential average of the well thickness of the casing metal still in place. In particular, this survey detects changes of the wall thickness and internal diameter of downhole casing. A typical inspection instrument for performing this survey is the casing inspection instrument marketed by Dresser Atlas under the name \"Magnelog tool\" and disclosed in the publication \"Dresser Atlas Casing Evaluation Services,\" pages 4-5 and 61-79; said publication being herein incorporated by reference. During operation, an inspection instrument generates an alternating (AC) magnetic field of low frequency. As the tool passes through the survey interval, the field permeates the casing wall and is detected by the receiver portion of the tool. The received electromagnetic wave is both attenuated and phase shifted by the casing wall, with the degree of phase shift in relationship to the wall thickness (weight). The degree of phase shift is detected and transmitted over the wireline for surface recording.\nThe phase shift curve registers the phase shift between generation of the low frequency electromagnetic wave and reception of the wave by the tool's receiver portion. The degree of phase shift increases with the increasing thickness of the casing's wall. The degree of phase shift may then be used in the calculation of the average thickness of the casing wall. While such a survey adequately calculates average thickness, it cannot detect the least wall thickness of casing or the degree of \"out-of-roundness\". (The term \"out-of-roundness\" generally refers to any non-circumferential internal wear, particularly one-sided wear) Thus, it is a distinct possibility that such a survey may fail to identify dangerously worn casing.\nA second survey available for assessing internal wear in downhole casing generates a depth recording of the greatest and least internal radii of the casing, thus allowing quantitative assessment of the degree to which the casing is internally worn or out-of-round. A typical instrument for performing this survey is the casing inspection instrument marketed by Dresser Atlas under the name \"Multi-Finger Caliper tool\" or \"MFC tool\" and disclosed in the publication \"Dresser Atlas Casing Evaluation Services\", pages 6 and 83-103; said publication having been previously incorporated by reference into this document. During operation, the MFC tool, which may include from 40 to 60 \"feeler arms\" which extend outwardly from the body of the tool in all directions, is centralized and all of the feeler arms extended outwardly from the tool for contact with the entire circumference of the inside casing wall at a particular depth. As the MFC tool moves through the casing, the feeler arms follow the contour of the casing's inside wall. The feeler arm with the maximum extension operates a first potentiometer (or other alternative device) which senses the greatest measured internal radius. The feeler arm with the least extension operates a second potentiometer (or other alternative device) which senses the smallest measured internal radius. The signals from the potentiometers are coupled through the electronics section and transmitted by the wireline to a recorder for presentation. While the MFC survey allows the evaluation of the degree to which the casing is internally out of round and permits estimating the greatest and least remaining wall thickness, the MFC survey cannot distinguish casing that is seriously worn on one side from unworn casing that is merely elliptical; in other words, severe one-sided wear may not be detected at all.\nThus, while both surveys are able to provide useful information regarding the internal wear of casing, both surveys fall short of being capable of providing a comprehensive evaluation of internal wear. For these reasons, a method for a combined casing evaluation survey became necessary so that an analytical technique which provided a comprehensive evaluation of internal casing wear would be possible. While separate evaluation of the results of the two separate survey followed by integrating the information derived from the surveys is possible, such a practice proved unworkable in actual use. A better method of achieving this goal of a comprehensive survey of internal casing wear is to devise a method by which selected data of the two previous surveys are combined and subject to a single analysis technique which will yield the desired more comprehensive internal casing wear evaluation."}
{"text": "An electroabsorption modulator (EAM) is a photonic semiconductor device that allows the intensity of a laser beam traveling through it to be controlled via an electric voltage. The principle of operation of the modulator is based on applying an electric field to cause a shift of the material absorption threshold toward longer wavelengths. A laser beam with a wavelength longer than that of the absorption threshold will either be transmitted without applied electric field or absorbed by applied electric field. A typical EAM has a waveguide structure and electrodes for applying an electric field in a direction that is perpendicular to the modulated light beam. In order to achieve a high extinction ratio, EAMs typically include a quantum well structure in the core of the waveguide. Due to the Quantum-Confined Stark Effect (QCSE), the quantum-well structure provides a strong shift of the absorption threshold at low electric fields. EAMs based on QCSE are capable of operating at relatively low voltages and at very high speeds (e.g., gigahertz (GHz)), which makes them suitable for use in optical communications systems.\nA distributed feedback laser (DFB) is a laser in which the entire cavity is made up of a periodic refractive-index structure that functions as a distributed reflector at the lasing wavelength, and contains a gain medium. An EAM can be monolithically integrated with a DFB on a single chip to form an electroabsorption modulated distributed feedback laser (EML) device that is capable of operating as a data transmitter. EML devices have been proposed that operate at high data rates (e.g., 10 gigabit per second (Gb/s)) in the 1550 nanometer (nm) and 1310 nm ranges. There are many difficulties that must be overcome in order to manufacture large volumes of these devices with high yield.\nFIG. 1 illustrates a cross-sectional top view of a known EML device 2 comprising a DFB 3 and an EAM 4. The DFB 3 and the EAM 4 each have metallic contact pads 5 and 6, respectively. One end facet 7 of the EML device 2 comprises a highly-reflective (HR) or anti-reflective (AR) coating. The other end facet 8 of the EML device 2 comprises an AR coating. An inter-contact isolation region 9 electrically isolates the DFB 3 and the EAM 4 from each other. The lines labeled 11 A and 11 B represent a straight reverse-mesa ridge structure 11 that extends across the DFB 3, the isolation region 9, and the EAM 4. The lines 13A and 13B delineate the channels located on either side of the ridge waveguide structure 11, which are created when the ridge waveguide structure 11 is formed.\nThe EML device 2 has trenches 12A and 12B formed therein to remove defect areas that exist after a surface area growth (SAG) process has been used to grow a multiple quantum well (MQW) structure, as described below in more detail with reference to FIG. 2A. The trenches 12A and 12B are generally uniform in width through the DFB 3, the EAM 4, and the isolation region 9. The trenches 12A and 12B are generally in parallel with and run alongside the reverse-mesa ridge structure 11.\nFIG. 2A illustrates a cross-sectional view of the EML device 2 shown in FIG. 1 along the A-A′ cross-section shown in FIG. 1. The different semiconductor layers that make up the EML device 2 are as follows. An n-type (001) Indium Phosphide (InP) material 21 is used as the substrate for the assembly 2. The substrate 21 has an n-type InP buffer layer 22 formed thereon. A MQW active region 23 is grown on top of the buffer layer 22 by a process known as Selective Area Growth (SAG). A p-type InP spacer layer 25 is grown on top of the MQW layer 23. A p-type Indium Gallium Arsenide Phosphide (InGaAsP) etch-stop layer 26 is grown on top of the spacer layer 25. Another p-type InP spacer layer 27 is grown on top of the InGaAsP etch-stop layer 26. A p-type InGaAsP grating layer 28 is grown on top of the spacer layer 27. A grating layer 28 is selectively etched in the DFB portion 3 to form a periodically varying refractive index region 31, which provides a filter for the laser spectrum in a desired wavelength value. A p-type InP infill 29 and cladding layer 32 are re-grown on top of layer 28. A p-type InGaAs contact layer 33 is then grown on top of the cladding layer 32.\nFIG. 2B illustrates a cross-sectional view of the DFB portion 3 of the EML device 2 along the E-E′ cross-section shown in FIG. 2A. After the contact layer 33 is grown, a silicon oxide (SiO2) dielectric mask 34 represented by mask portions 34A, 34B and 34C is deposited on the top of the contact layer 33 and a wet chemical etch process is performed to etch the contact layer 33 and the cladding and infill layers 32 and 29.\nFIG. 3 illustrates an enlarged view of the DFB portion 3 of the EML device 2 along the E-E′ cross-section shown in FIG. 2A after the chemical etching process has been performed. When the chemical etch process is performed, the device 2 etches relatively quickly in the (001) plane, but relatively slowly in the (111) plane, causing the waveguide ridge 35 to be formed with lateral (111) facets 36 thereon, having an angle, ω, between the (001) and (111) facets that is typically about 54.7 degrees. This ridge configuration is commonly referred to as a reverse-mesa ridge configuration due to the fact that the top of the ridge is wider than the base of the ridge.\nFIG. 4A illustrates a cross-sectional view of the DFB 3 shown in FIG. 2A along the E-E′ cross-section prior to the layers of the MQW active region being grown using the SAG process. Prior to performing the SAG process to grow the MQW active region, a dielectric mask pattern comprising a plurality of rectangular mask portions 54A and 54B are deposited on top of the substrate 21, leaving openings 22 in which to grow the MQW active region.\nFIG. 4B illustrates a cross-sectional view of the DFB 3 along the E-E′ cross-section shown in FIG. 2A after the layers making up the MQW active region 23 have been grown. The layers 23 are grown only on the areas of the surface of the substrate 21 that are exposed through the mask pattern 54A and 54B.\nFIG. 4C illustrates a cross-sectional view of the EAM 2 along the E-E′ cross-section shown in FIG. 2A after the mask pattern 54A and 54B has been removed. After the mask pattern 54 has been removed, it can be seen that the surface is essentially non-planar. This non-planar surface is difficult to work with in performing the additional processes described above with reference to FIGS. 2A and 2B, e.g., growing the p-type InP spacer layer 25 on top of the MQW layer 23, growing the p-type InGaAsP etch-stop layer 26 on top of the spacer layer 25, growing the p-type InP spacer layer 27 on top of the InGaAsP etch-stop layer 26, growing the p-type InGaAsP grating layer 28 on top of the spacer layer 27, forming the periodically varying refractive index region 31 in the grating layer 28, etc., up through and including the formation of the reverse-mesa ridge structure 35. Because of the difficulties of working with the non-planar surface that results after the SAG process has been performed, the resulting assemblies may have performance problems, which will typically result in reduced manufacturing yield and increased costs.\nWith reference again to FIG. 1, it is important for the isolation region 9 between the DFB 3 and the EAM 4 to have high resistance and low capacitance. One technique that is used to provide the isolation region 9 with high resistance is to perform ion implantation in order to implant ions in the isolation region 9 to increase its resistance. One disadvantage of this technique is that it increases process complexity, which increases costs. Another disadvantage of this technique is that it is not always effective at providing the isolation region 9 with a sufficiently high resistance, which can lead to performance problems, and consequently, to reduced yield and increased costs.\nA technique that is sometimes used to provide the isolation region 9 with decreased capacitance involves depositing a thick dielectric material underneath the contact pads 5 and 6. One disadvantage of this technique is that it increases process complexity, which increases costs. Another disadvantage of this technique is that it sometimes results in the contact pads 5 and 6 failing to adhere to the EML device 2. In other words, the pads 5 and 6 can peel away from the assembly 2, resulting in performance problems, and consequently, in reduced yield and increased costs.\nAccordingly, a need exists for a photonic semiconductor device, such as an EML device of the type described above, for example, that is capable of operating at high speeds (e.g., 10 GB/s) over a wide temperature range in the 1310 and 1550 nm windows. A need also exists for a method for manufacturing such a device in such a way that manufacturing yield is kept relatively high and costs are kept relatively low."}
{"text": "Electromagnetically actuated switches exist in a variety of forms and for a variety of functions. Typically, they serve the function of opening and closing electrical circuits by controlling the bridging of a pair of contacts by a movable conductor. That conductor forms part of, or is carried by, an armature, the displacement of which is controlled by the selective energization of an electrical coil. In most such instances, the coil remains energized to maintain the contacts in one of their operative states, either open or closed. However, in many other instances, it may be necessary or at least desirable to maintain the associated switching circuitry in each of its stable states without requiring continued energization of the coil.\nExemplary of one type of the latter mentioned class of electromagnetic switches are U.S. Pat. Nos. 820,119; 1,824,973; 2,528,520; 2,851,646; 2,892,059; and 3,401,892. Each of these patents depicts an electromagnetic switch in which a linearly-moving armature or plunger is caused to rotate as a result of linear reciprocation. Such rotation is operative to move a bridging conductor between successive, angularly-displaced contacts, typically for some form of sequential control.\nExemplary of another type of electromagnetic switches of the class described earlier are U.S. Pat. Nos. 1,908,567; 2,703,348; 2,874,244; and 2,922,861. These electromagnetic switches are typically concerned with simply moving the contact-bridging element between either one of two stable states, e.g. \"Contacts Closed\" or \"Contacts Open\". Apart from the brief application of an actuating potential to the coil for moving the contact-bridging structure between its two stable states, no further energization of the coil is required. Such switches have, in some instances, been characterized as bistable relays or solenoids and will be so characterized herein. Typically, these electromagnetic switches are of a solenoid configuration in which an armature or plunger is reciprocated axially of an annular coil. Moreover, some portion of the plunger is usually axially coextensive with the coil during some portion of operation. In such bistable solenoid switches, it has been relatively common to effect the requisite bistable operation by means of a pin-and-slot arrangement. A pin, carried by the plunger, is subjected to the camming action of a stationary slot during axial displacement to provide the requisite rotational control of the plunger. U.S. Pat. Nos. 1,908,567; 2,874,244; and 2,922,861 are of that type. In another configuration represented by the aforementioned U.S. Pat. No. 2,703,348, the pin may be stationary and the slot may be formed in the plunger. The pin extends radially inward from the coil and into the slot.\nThe provision of such pin-and-slot arrangement to obtain a bistable solenoid switch complicates the manufacturing process and thereby increases its cost. Moreover, the integrity of the resulting system may, in some respects, be compromised. Certainly the assembly of the pin-and-slot structure in the aforementioned U.S. Pat. No. 2,703,348 patent presents a significant challenge."}
{"text": "A communication system is formed, at a minimum, by a transmitter and a receiver interconnected by a communication channel. The transmitter forms at least a portion of a sending station, and the receiver forms at least a portion of a receiving station. Communication signals transmitted by the sending station are transmitted upon the communication channel to be received by the receiving station. Information contained in the communication signals transmitted by the sending station is recovered once received at the receiving station. In a digital communication system, information to be communicated to the receiving station is digitized. The digitized information is then used to form the communication signal.\nIn some conventional, digital communication systems, the communication signals are transmitted in bursts. Blocks of digital information are communicated by a sending station to a receiving station during transmission of the bursts upon the communication channel.\nIn a non-ideal communication system, the communication signal is distorted during its transmission upon the communication channel. Because of such distortion, when the communication signal is received at the receiving station, the received signal differs somewhat from the communication signal transmitted from the sending station. If the communication channel is of poor quality, and the amount of distortion is significant, the informational content of the communication signal cannot be properly recovered at the receiving station.\nIn a digital, radio communication system, for instance, multi-path and Rayleigh distortion is sometimes introduced upon the communication signal as it is transmitted on the communication channel from the sending station to the receiving station.\nVarious schemes have been developed to better ensure that the informational content of the communication signal transmitted upon a non-ideal channel can be recovered at a receiving station. Several of such schemes utilize a feedback arrangement in which the receiving station reports back to the sending station whether the informational content of the received signal can be adequately recovered.\nSeveral of such schemes have been implemented in a digital communication system in which blocks of information are communicated in bursts. The receiving station determines whether a burst of the block of information received at the receiving station is of an acceptable quality level. Indications of the determinations, forming feedback information, are returned to the sending station by way of a feedback channel. The indication may also be communicated back to the sending station implicitly. That is to say, lack of transmission of feedback information from the receiving to the sending station may be indicative of successful recovery of the informational content of the block of digital information. Or, the feedback information may be provided to the sending station in a manner other than by way of a direct feedback channel extending directly to the sending station. The indication may alternately be communicated back to the sending station indirectly. For example, the sending station might measure the channel conditions and decide that the probability of the receiving station successfully decoding the block of digital information to be so low that retransmission of the block of digital information would be justifiable.\nMost simply, the receiving station simply detects whether the received signal burst is of an acceptable quality level. If the signal quality level of the received signal burst is not good enough to permit recovery of the informational content thereof, the receiving station merely requests the sending station to retransmit the block of digital information in a subsequent burst. Such a scheme is sometimes referred to as an ARQ (automatic request) scheme.\nRetransmission of the block of digital information can be repeated in successive bursts until the block of digital information is received at the receiving station with at least a minimum quality level. The receiving station makes the determination of the quality of the received block of digital information responsive, for example, to detection of an error detecting code, knowledge of the communication channel upon which the block of information is transmitted, or some other suitable scheme.\nThe block of digital information can also be transmitted using an error protecting code, such as type I hybrid ARQ (Automatic Repeat Request) scheme, a general channel decoding technique. When the receiving station receives the block of digital information, the received block is decoded by a decoder to extract the informational content from the received signal. A determination is made if the informational content of the block of digital information can be recovered with at least an acceptable quality level. Because of the error protection, the block of digital information is better able to be recovered even if transmitted upon a communication channel of lowered quality.\nFurthermore, if the decoder has the ability to exploit, not only the values of the received symbols, but also reliability information on the symbols, so-called soft information, the performance will increase substantially. Decoder circuitry at the receiving station capable of utilizing soft information is referred to as a soft input decoder.\nIn some other schemes, such as a type II and type III, hybrid ARQ (Automatic Repeat Request) scheme, blocks of digital information determined by the receiving station to be corrupt, i.e., to be of poor quality, are not merely discarded. Instead, the corrupt blocks are combined with subsequently-transmitted blocks of digital information. Information contained in previously-transmitted blocks is accumulated and forms accumulated knowledge. Such accumulated knowledge can be used to facilitate recovery of subsequently-transmitted blocks. Thereby, \"accumulated knowledge\" results, the number of times in which the blocks of digital information must otherwise be retransmitted decreases. The advantages of utilization of accumulated knowledge permitted by the combining of successive retransmissions of the blocks of digital information occur with increased amounts of information contained in the blocks of digital information.\nThis is advantageous as the probability of correct decoding generally increases if the accumulated knowledge is used during the decoding rather than if only the latest retransmitted block is utilized. Decoder circuitry at the receiving station capable of utilizing soft information is referred to as a soft input decoder.\nIn another scheme, a sequence formed of other symbols created from the same block of digital information is instead transmitted. For instance, if the first sequence is formed of an original set of parity symbols (possibly including uncoded information symbols of the block of digital information) and retransmission is requested, additional numbers of parity symbols are transmitted by the sending station in the retransmitted sequence. The receiving station accumulates the received symbols of the various transmissions and the symbols accumulated during the various transmissions are utilized together during a decoding process.\nA conventional block or convolutional low rate coding process can be performed to form an encoded signal. Then, the encoded signal is \"punctured.\" That is, selected symbols of the coded block are marked so as not to be transmitted. Only the \"unpunctured\" symbols, i.e., symbols which have not been marked, are transmitted. While the coded signal is weakened by transmitting only some of the encoded symbols, a higher effective coding rate is possible for a given encoder structure. If the informational content of the block of digital information can not be recovered with a desired level of quality from the symbols first received at the receiving station, additional coded portions of the signal are requested by the receiving station to be transmitted thereto. Some of the previously-punctured, and not-yet sent, symbols are thereafter transmitted by the sending to the receiving station. The receiving station utilizes both the previously-transmitted symbols and the newly-transmitted symbols. If additional retransmissions are required, yet additional symbols are subsequently transmitted to the receiving station. Such a scheme is referred to as a successive redundancy transmission scheme.\nA successive redundancy transmission scheme, designed to be able to cope with adverse channel conditions, however, requires a complex decoding process to be performed to recover the informational content of the block of digital information. And, if the block of digital information is coded using a convolutional coding scheme, the decoding required at a receiving station to decode the information is complex irrespective of the amount of puncturing of the code.\nWhen using a punctured convolutional code, a set of received redundant symbols that would, by themselves, cause a decoding error to occur still cooperate in the same way to counteract any error correcting attempts supported by additional redundant bits.\nWhen the channel conditions of the communication channel extending between the sending and receiving stations are of poor quality, conventional codes of moderate complexity do not perform well. High-complexity codes, which perform well at such conditions, however, add unnecessary complexity when the channel conditions are of good quality. That is to say, while complex coding is necessitated when the channel conditions are of poor quality levels, such channel coding is not necessary when the channel conditions are of good quality levels.\nA manner by which to adaptively select the complexity of the coding of a block of digital information to be transmitted between a sending and a receiving station would therefore be advantageous. When channel conditions are of good quality, only limited portions of a coded signal would have to be transmitted to a receiving station to permit recovery of the informational content thereof. And, when channel conditions are of poor quality, additional portions of the coded signal could be transmitted, better to ensure that the informational content of the block of digital information can be recovered at the receiving station.\nIt is in light of this background information related to digital communication systems that the improvements of the present invention have evolved."}
{"text": "The ocean has great potential for generating usable energy if it can be harnessed efficiently. For example, ocean waves, high and low ocean tides, and/or temperature differences in the water are a several ways that the ocean can be used to generate useable energy. Ocean waves, in particular, can have a significant amount of kinetic energy and this energy can be used to power various systems. Although there are many systems for generating energy from the movement of ocean water, there is a continued need for improvements in transferring energy from a power generation system to another location, such as a power grid on shore. Some such devices utilize a buoyant tube to transfer such energy."}
{"text": "This invention relates to preforms for blow-molded articles, and more particularly, to antinesting preforms for blow-molded synthetic resinous containers.\nFood and beverage containers formed of PET (polyethylene terephthalate), and particularly PET bottles for carbonated or effervescent beverages, have become ubiquitous throughout the principal industrial countries of the world.\nPET bottles are commonly formed by the reheat-stretch-and--blow procedure. Though this procedure is well known, it may be described briefly as follows.\nTubular preforms, closed at one end and open at the other, are injection-molded of a suitable grade or formulation of PET. The neck finish, that is, the finish at the open end, is formed in the injection mold in a configuration to receive a suitable closure when the resulting container is filled.\nThe preforms are removed from the injection mold and allowed to cool to ambient temperature, and are later introduced to a blow-molding system where they are first heated to a temperature within the glass-transition range of the material, then mechanically stretched in the longitudinal direction, and finally blow-molded to final container configuration and dimensions. The mechanical stretching and blow-molding impart biaxial molecular orientation to the material of the container body, thereby enhancing the tensile strength of the body wall while reducing its thickness.\nA number of machines and systems for performing this procedure are well known and commercially available.\nThe finished containers take a great variety of sizes and configurations. The problem of designing a PET container for a particular purpose presents a bewildering variety of considerations, among them not only capacity and aesthetic appeal, but economy of material, and importantly, whether the container is to be filled with contents under pressure, as in the case of carbonated beverages, or under vacuum, as in the case of hot-filled products.\nWhen container configuration has been determined, an appropriate preform must be designed. Because of the need to conserve material and yet ensure that the material will be so distributed in the finished container that it will exhibit adequate resistance to deformation and rupture under all conditions of expected use, preform design has become an exact science, and an exacting one.\nIn the case of certain PET bottles intended to package highly carbonated beverages subject to internal pressures of three to four atmospheres, a preform configuration has evolved in which the outer diameter of the closed end is smaller than the internal diameter of the open end, whereby one preform may, by way of its open end, receive the closed end of an adjacent preform. Simply stated, such preforms are subject to nesting with one another.\nNesting is an impediment to safe and efficient handling. Heretofore, it was necessary to detect and separate nested preforms before damage could occur to them or to the equipment in which they were being processed. In many cases, detection and separation depended exclusively upon the watchfulness and dexterity of a human operator."}
{"text": "In the art it is desirable to produce elastomeric compounds exhibiting reduced hysteresis when properly compounded with other ingredients such as reinforcing agents, followed by vulcanization. Such elastomers, when compounded, fabricated and vulcanized into components for constructing articles such as tires, power belts, and the like, will manifest properties of increased rebound, decreased rolling resistance and less heat-build up when subjected to mechanical stress during normal use.\nThe hysteresis of an elastomer refers to the difference between the energy applied to deform an article made from the elastomer and the energy released as the elastomer returns to its initial, undeformed state. In pneumatic tires, lowered hysteresis properties are associated with reduced rolling resistance and heat build-up during operation of the tire. These properties, in turn, result in lowered fuel consumption of vehicles using such tires.\nIn such contexts, the property of lowered hysteresis of compounded, vulcanizable elastomer compositions is particularly significant. Examples of such compounded elastomer systems are known to the art and are comprised of at least one elastomer (that is, a natural or synthetic polymer exhibiting elastomeric properties, such as a rubber), a reinforcing filler agent (such as finely divided carbon black, thermal black, or mineral fillers such as clay and the like) and a vulcanizing system such as sulfur-containing vulcanizing (that is, curing) system.\nPrevious attempts at preparing readily processable, vulcanizable, silica-filled rubber stocks containing natural rubber or diene polymer and copolymer elastomers have focused upon the sequence of adding ingredients during mixing (Bomal, et al., Influence of Mixing procedures on the Properties of a Silica Reinforced Agricultural Tire Tread, May 1992), the addition of de-agglomeration agents such as zinc methacrylate and zinc octoate, or SBR-silica coupling agents such as mercapto propyl trimethoxy silane (Hewitt, Processing Technology of Silica Reinforced SBR, Elastomerics, pp 33-37, March 1981), and the use of bis[3-(triethoxysilyl)propyl]tetrasulfide (Si69) processing aid (Degussa, PPG).\nThe use of Si69 processing aid in the formulation of silica-filled rubber stocks has been successful, but generally requires a large amount of the additive, such as 10% by weight based on the weight of silica, in order to be effective.\nPrecipitated silica has been increasingly used as a reinforcing particulate filler in carbon black-filled rubber components of tires and mechanical goods. Silica-loaded rubber stocks, however, exhibit relatively poor processability.\nThe present invention provides a mixture of silanes for use as processing aids for silica-filled rubber stocks, which greatly improve the processability and properties of the formulations and resulting vulcanized product."}
{"text": "The field of art to which the invention relates includes motors and pumps, and more especially fixed displacement fluid motors and pumps.\nIt has long been a problem in the art to which this invention pertains to effectively disengage during certain modes of operation of fluid motors and pumps of the type contemplated, the cam track and pistons such that during high speed free-wheeling operation, for example, of the motor or pump the pistons are held in inactive retracted positions without engaging the driving cam surface at any point in the rotary cycle of operation, thereby both protecting the motor or pump from damage as a result of the pumping action of the pistons at excessive rpm, and minimizing the wear and tear of parts associated with reciprocating parts. Heretofore, to my knowledge, no fully satisfactory means has been devised for effectively operating at high speed, fixed displacement fluid motors or pumps while holding the pistons thereof in fixed inoperative positions.\nIn the use of hydrostatic wheel motors in certain types of vehicles, such as an off-the-road construction type tractor vehicles, for example, it is often necessary to tow the vehicle on highways from one work location to another. Heretofore in certain applications highway towing speeds have been relatively low in order to minimize wear and tear on parts such as pistons, cylinders and bearings, and so that the motors were not damaged or destroyed by high speed towing operation beyond the design limits of the motors.\nAgain, in operating certain types of cable winches it is desirable, for example, to disengage the winch drive from time to time in a high speed cable pull-off operation. My invention is adapted to effectively disengage such a winch drive motor to stop the piston action which would otherwise occur. Also, for example, it is frequently desirable to accelerate a fly wheel to a relatively high speed to store energy, subsequent to which the fly wheel speed and energy is conserved by minimizing losses. During such fly wheel acceleration powered by the engaged drive motors it is then desirable to disengage the pistons by holding them in retracted positions, and following fly wheel acceleration to conserve fly wheel speed and energy and then extract the stored energy as needed from time to time by reenergizing the pistons when drawing power from the flywheel.\nFixed displacement camming ball or roller radial or axial piston type air compressor motors, steam engines and pumps may be used in a variety of applications in which the kind of problem indicated above exits. My invention may be applied to all such exemplary types of fluid motors or pumps to provide overspeed protection, and/or reduce wear and tear on the motor or pump under certain operational modes, and/or minimize the noise and heat losses generated in such modes of operation."}
{"text": "1. Field of the Invention\nThe present invention is directed toward a zipper attachment device, and more particularly to such devices that permit a user to selectively suspend, remove and interchange differing decorative designs, logos or jewelry from the apertures of zippers to both add to the operation of the zipper and add to its aesthetic value.\n2, Prior Art Statement\nOver the years, zipper apertures have played host to a variety of suspended articles. Traditionally such devices were made of cloth or leather, to assist in the operation of the zipper, or of string or plastic, to suspend informational information about the zippered item. The traditional cloth and leather attachments are tied to zippers, through their apertures, and extend a length downward to assist a user in gripping and operating the zipper. Such extensions may have items attached to their ends, but such attachments are usually sewn or tied to the end and are not easily exchanged or removed.\nTraditional informational zipper attachments are usually made from a thin strand of plastic or string and serve as mediums to suspend printed information such as price tags, warranties, manufacturer identification and the like. These attachments have no functional value to the zipper and are suspended from the zipper solely as not to damage the material of the zippered item. As such, the informational attachments on zippers are made to be easily removed, and once removed not reattachable to the zipper aperture.\nThere exists no prior art concerning devices, such as the present invention, for selectively suspending decorative designs and logos from zipper apertures. The present invention can suspend a design from a zipper and lock it into place, to assist in the operation of the zipper, however the invention allows the suspended design to be repeatedly adjusted, removed and replaced without harming the suspended design or the suspension medium. The changing of the suspended design is quick and easy and does not require untying or other complicated manipulations of traditional means. Additionally, the present invention provides for an attachment of a zipper extension that not only assists in the operation of the zipper, but lets the operator change the aesthetics of the zipper extension at will.\nThus, prior art does show the use of zipper apertures to suspend functional or informational items, but prior art does not show a zipper suspension device that allows for the quick and easy interchanging of decorative designs, that are reusable, adjustable, have aesthetic value and serve to functionally assist a zipper's operation."}
{"text": "Sensors convert nonelectrical physical or chemical quantities into electrical signals. The behavior of semiconductor sensor devices is modeled, under a series of assumptions, by a system of nonlinear partial differential equations and associated boundary and input value conditions. For such sensors, it is essential to solve numerically at least some of the equations in order to obtain a meaningful result describing environmental conditions encountered by the sensor. Mechanical sensors also suffer from some of the same issues relating to nonlinearity and associated boundary and initial value considerations.\nWhen the desired result is determined by the integration of sensor output over time, errors of bias scale, and nonlinearity occur. Mathematical algorithms are used to improve the accuracy of the result by correcting for such errors of bias, scale and nonlinearity.\nThe present invention relates to methods and devices of determination of an input signal which changes in time as well as of its integral value.\nMethods are known of the above mentioned general type and disclosed for example in Inventor's Certificates of the USSR 1,453,418 and 1,541,635. In accordance with the method and device disclosed in these references, integration is performed to a given value of the integration result. This objective is presented for example in inertia systems of targeting, for turning off an engine when a rocket reaches a desired speed V.sub.r. The input signal in this case is an acceleration which is measured by accelerometer, and its inaccuracies are a main source of error of the system as a whole. The device includes a convertor 1 for a frequency of pulse sequence, a frequency multiplier 2, an integrator or pulse counter 3, a source of reference signals 4, keys 5 and a control device 7. A code corresponding to a desired value of speed is introduced into the pulse counter by integration of reference signals (accelerations) during a time corresponding time interval. The values of the reference signals can be obtained by turning a sensitivity axis of the accelerometer relative to a local vertical. Then, during measurements of the input signal which is the acceleration during a flight, an inverse integration is performed. In other words, the pulses from the convertor, multiplied many times by the frequency multiplier are deducted in the counter from the code obtained during the time of direct integration when the counter was operating for addition of the reference signals. A system of equations was used to determine the reference signals and time intervals. The nature of change of the input signal was approximately known beforehand, and therefore an initial moment of the input signal was calculated. Zeroing of the pulse counter is a signal that the speed reached the desired value. The above described method has a high accuracy. However, it has some disadvantages, namely the fact that the high accuracy is guaranteed in only one point and not along the whole scale of speed changes. This values must be known in order to determine of the location of a rocket. The known method and device are not always technically implementable since it is necessary that one reference signal is greater and the other is smaller than an average value of the input signal. However, the acceleration can reach a few g while on the ground during giving of the reference signals not always there is a standard of acceleration more than 1 g."}
{"text": "1. Field of the Invention\nThe present invention is directed to file folders, and more particularly to stair-step style expandable file folders.\n2. Description of the Related Art\nFile folders are common in home and office settings.\nA typical expandable file folder includes a series of rectangular dividers that are joined on three sides to respective pleats in an expandable gusset. The folder has pockets formed between each adjacent pair of dividers. Each pocket is accessible from the unbound edge of the dividers, which will be referred to in this document as the top edge of the file folder.\nOne common type of expandable file folder is referred to as a stair-step expandable folder. In a stair-step expandable folder, the dividers extend by different amounts along the top edge of the folder. The dividers at the back of the folder extend upward farther than those at the front of the folder, so that the top edges of all the dividers are viewable simultaneously. A user can select a desired pocket directly, rather than having to flip through all the dividers to find the correct pocket.\nOne potential drawback to typical stair-step expandable folders is that the pockets have different depths. Because the dividers at the rear of the folder extend farther upward than those at the front of the folder, the pockets at the rear of the folder are deeper than those at the front of the folder. As a result, documents stored in a rear pocket are more difficult to access than those stored in a front pocket. This difficulty in accessing the rear pockets is unacceptable.\nAccordingly, there exists a need for a stair-step expandable file folder in which the stored documents are easily accessible for all pockets."}
{"text": "Efficient list decoding beyond half a minimum distance for Reed-Solomon and Bose, Ray-Chaudhuri and Hocquenghem (i.e., BCH) codes were first devised in 1997 and later improved almost three decades after the inauguration of an efficient hard-decision decoding method. In particular, for a given Reed-Solomon code C(n,k,d), a Guruswami-Sudan decoding method corrects up to n−√{square root over (n(n−d))} errors, which effectively achieves a Johnson bound, a general lower bound on the number of errors to be corrected under a polynomial time for any code. Schmidt, Sidorenko, and Bossert devised a multi-sequence shift-register synthesis to find an error locator polynomial beyond half a minimum distance for low-rate (i.e., <⅓) Reed-Solomon codes when a unique solution exists. Apart from small probability of failure due to ambiguity, the resulting decoding radius is identical to that of the Sudan technique. The Sudan technique extended the Guruswami-Sudan technique to achieve subfield Johnson bounds for subfield subcodes of Reed-Solomon codes by distributing multiplicities across the entire subfield.\nWu presented a list decoding technique for Reed-Solomon and binary BCH codes. The Wu list decoding technique casts a list decoding as a rational curve fitting problem utilizing polynomials constructed by a Berlekamp-Massey technique. The Wu technique achieves the Johnson bound for both Reed-Solomon and binary BCH codes. Beelen and Hoeholdt re-interpreted the Wu list decoding technique in terms of Gröbner bases and an extended Euclidean technique to list decoded binary Goppa codes up to the binary Johnson bound.\nIt would be desirable to implement a combined Wu and Chase decoding of cyclic codes."}
{"text": "A luminaire (i.e., a light fixture) is typically a complete lighting unit consisting of one or more lamps, socket, and optical devices for distributing light. A recessed luminaire is an example of a luminaire that is installed behind a structure such as a ceiling. Recessed luminaires are used in both commercial and residential applications. For example, a recessed luminaire is typically mounted on support structures behind a ceiling wall that has an opening to allow light from the recessed luminaire to illuminate an area below the ceiling wall.\nVarious support systems have been employed to support recessed luminaires. For example, recessed luminaires are often suspended between support structures such as joists and T-bar structures. To illustrate, a recessed luminaire may be supported by hanger bars that extend between parallel support structures. Generally, hanger bars need to have adequate strength and rigidity to reliably support a recessed luminaire. Further, because the spacing between the support structures such as the joists of a ceiling structure may vary, adjustability of the lengths of hanger bars is desirable for easy installation as well as for compatibility with different support structures. However, unintended detachment of hanger bar members from each other can make installation of hanger bars difficult, which can waste time and money.\nThus, adjustable hanger bars that have structures that prevent or reduce risk of unintended detachment of the hanger bars from each other are desirable."}
{"text": "For preparing catalysts, there have been various methods of reducing catalyst precursors containing metal oxides proposed. For example, a copper-based catalyst is used in production of an alcohol by hydrogenation of a carboxylic acid or a carboxylic acid ester. When a fixed-bed reaction system is employed, the copper-based catalyst is entirely subjected to gas phase reduction for activation. In industry, such gas phase reduction of a catalyst is generally carefully conducted at a predetermined temperature under an inert gas flow that contains several to several tens percent of hydrogen in order to avoid local overheat due to rapid reduction of the catalyst.\nIn some cases, reduction of a metal oxide with hydrogen generates a large amount of heat. For example, one mole of copper oxide generates 20 kcal of reduction heat. In addition, reduced copper is known to have very low thermal stability. Thus, it is important to gradually conduct reduction with controlling heat generation to produce a copper catalyst without impairing its performance. This is particularly important for a molded catalyst, because it has a difficulty of heat release.\nTherefore, it is naturally expected for a catalyst subjected to reductive activation with high concentrate hydrogen for a short time in a gas phase to have significantly decreased catalytic performance due to rapid generation of heat, and more for a catalyst subjected to reductive activation for a short time in such a large scale as in industry to fall into very serious situations due to rapid elevation of temperature. As thus, a general method for practical gas phase reductive activation of a catalyst containing copper oxide is generally conducted with low concentrate hydrogen for a long time. For instance, JP-A 61-161146 describes that such reductive activation takes a time from 4 to 14 days.\nAs described above, fixed-bed reaction systems generally employ gas phase reduction, but some systems employ liquid phase reduction to activate a catalyst precursor containing copper oxide. For example, JP-B 2990568 discloses liquid phase reduction of a molded precursor of a copper-containing hydrogenation catalyst, that reduction is conducted at a temperature ranging from 50 to 140° C.\nJP-B 3195357 discloses activation of a copper-containing hydrogenation catalyst by supplying a mixture of hydrogen with an inert gas under an inert solvent flow."}
{"text": "It has been known in the plastic pipe industry for several decades that pipe couplings made by injection molding are inherently weaker than couplings made from extruded material. Further, accelerated life tests and experience with old pipe lines shows that these couplings have a life substantially less than that of extruded pipe which they join together. Part of the problem is due to the strength of the plastic materials. Plastic formulations for injection moldings are inherently weaker than the formulations that can be extruded. For example, polyvinyl chloride (PVC) that is extruded can easily have a design strength rating of 2,000 psi. However, to provide the necessary plasticity and flow characteristics for injection molding, it is difficult to use a PVC formulation that exceeds a design strength 1600 psi (design stress) for conventional injection molding machines.\nAnother part of the problem has to do with the injection molding process. The cavity of the mold is usually fed by a side inlet or gate, and the material flows around the inner core of the mold on each side of the gate. The liquid plastic being injected into the mold joins together approximately opposite the gate and more or less welds together along a line known as the weld line. This weld line can be a line of weakness where injection molded couplings fail upon being subjected to high internal pressures caused by the fluids in the pipe line. The various conditions of injection molding must be critically met in order to obtain a good weld line.\nAs a result, couplings have been largely replaced by bells or enlarged sockets formed on one end of each length of plastic pipe. A typical installation of a pipe line involves inserting the spigot or plain end of another section of pipe into an adjoining bell and the joint is sealed by a ring of rubber or rubber-like material held in an internal groove in the bell. The bells may be formed by properly expanding a pipe and thus the inherent strength of the extruded material and the extrusion process is preserved.\nHowever, these enlarged bells give rise to difficulties in manufacturing, storing and shipping plastic pipe. In stacking the pipe, each adjoining length has to be reversed, and each layer has to alternate plain ends and bell ends. The stacks are thus longer than the pipe sections and the protruding bells receive much damage in handling and in transportation. Furthermore, each length of extruded pipe has to be subjected to the belling process which adds substantially to the cost. For these and other reasons the plastic pipe industry has long sought a satisfactory and economical coupling.\nVarious attempts have been made to provide extruded couplings. Oversize extruded pipe sections have been internally machined to form the seal grooves and tapers, but these have proved to be expensive. Also, oversize extruded pipe sections have been internally expanded to create seal grooves in each end, but because of thinning at the grooves, the material had to be grossly oversize to start with, resulting in unnecessary material and cost."}
{"text": "1. Field of the Invention\nThe invention relates to an excimer laser pumped by an electrical gas discharge, and particularly to a preionization device and technique for generating a stable pulsed gas discharge for pumping of an active medium of an excimer laser.\n2. Discussion of the Related Art\nUV-preionization of the electrical discharge in a pulsed gas laser is typically realized by means of an array of spark gaps or by another source of UV-radiation (surface, barrier or corona gas discharges), disposed in the vicinity of at least one of the solid electrodes of the main discharge of the laser. Conventional pulsed electrical gas discharges typically used for pumping the active media of excimer lasers are unstable. The development of discharge instabilities cause the glow discharge, a precondition for laser emission, to have a short phase (e.g., having a typical duration from 10-100 ns) and to thus be terminated too quickly. The desired way of generating a high quality gas discharge for use in excimer lasers is to provide an intense, yet uniform preionization of the gas volume before the main gas discharge occurs.\nOne way of providing this preionization is by photo-ionizing the laser gas with UV-light emitted from an auxiliary gas discharge before the main gas discharge is switched on. Some known methods of preionizing high pressure gas lasers include x-ray, spark and corona-gap preionization. See R. S. Taylor and K. E. Leopold, Pre-preionization of a Long Optical Pulse Magnetic-Spiker Sustainer XeCl Laser, Rev. Sci. Instum. 65 (12), (December 1994). The spark method involves the use of spark gaps (ordinary or stabilized by a dielectric surface), and the corona-gap method involves the use of pulsed corona-like discharges near a dielectric surface.\nAreas of focus for design improvement of corona-gap preionizers include the geometry of the dielectric body, and the arrangement of the preionization electrodes. See U.S. Pat. No. 4,718,072 to Marchetti et al. (showing a grounded internal preionization electrode surrounded by a dielectric having a positive potential applied to its outer surface through contact with the positively biased main electrode), European Patent Application published) EP 0 532 751 A1 (showing an internal preionization electrode surrounded by a dielectric buried in one of the main electrodes); U.S. Pat. No. 4,953,174 to Eldridge et al. (showing the dielectric surrounding an internal preionization electrode abutting with a main discharge electrode); see also R. Marchetti et al., A New Type of Corona-Discharge Preionization Source for Gas Lasers, J. Appl. Phys. 56 (11), (Dec. 1, 1984); U.S. Pat. No. 4,380,079 to Cohn et al.\nReconfiguration of external electrical circuits is another area where corona-gap pre-ionizer design improvement efforts have been focused. See Taylor et al., citation above; U.S. Pat. No. 5,247,531 to Muller-Horsche (showing an excitation of preionization electrodes affected by the same high voltage source as the main discharge electrodes including a time delay inductance between them), U.S. Pat. No. 5,247,534 to Muller-Horsche (including flow bodies configured to facilitate laser gas flow and formed of material exhibiting secondary x-ray emission characteristics) and U.S. Pat. No. 5,247,535 to Muller-Horsche (disclosing electron emission from a heated cathode, wherein x-rays emitted as the electrons impinge upon a separate anode serve to preionize the laser gas volume).\nU.S. Pat. No. 5,337,330 to Largon describes the typical corona-like preionization arrangement of FIG. 1a. See also U.S. Pat. No. 5,247,391 to Qormley, and U.S. Pat. No. 4,953,174 to Eldridge et al. The discharge chamber having the preionizing device of FIG. 1a includes a high voltage main electrode 1 and a grounded main electrode 2. The preionization unit includes two internal preionization electrodes 3a each located on one side of main discharge region 5 between the main discharge electrodes 1,2. Each preionization unit includes a dielectric tube 3b of generally cylindrical shape surrounding the internal preionization electrode 3a. A preionization discharge (ultraviolet emission) 4 from the preionization electrodes 3a and 6 and dielectric tubes 3b causes a preionization of the volume of the main gas discharge. A pair of external preionization electrodes 6 of the preionization units comprise metal plates and are each directly connected to the nearby main discharge electrode 1 (e.g., the cathode at high potential).\nIn this case energy stored in the dielectric tubes 3b, which can be non-negligible relative to the energy of the main discharge, during a preionization phase, will also be absorbed into the main discharge 5. However, that added energy typically will not increase the laser output due to a high wave impedance of the dielectric tubes 3b. The tubes 3b act much like a charged transmission line in that this wave impedance is typically much higher than the impedance of the main gas discharge. The high wave impedance is caused by a distributed inductivity of the whole dielectric tubes 3b (as a transmission line) and a concentrated inductivity at the point of electrical connection of the tubes 3b with the internal corona discharge electrodes 3a. \nThe residual energy produces high voltage electrical oscillations between the capacitance of the dielectric tubes 3b of the preionization units and the main gas discharge volume. These high voltage oscillations are undesirable because they significantly reduce the ability of the dielectric tube 3b of the preionization unit to resist direct high voltage breakdown and over-flashing near the open ends of the dielectric tubes 3b. Moreover, these oscillations deteriorate the quality of the main gas discharge 5 and thus hinder the operation of the laser, particularly during operation at a high repetition rate. Furthermore, the oscillations cause additional wear to the main gas discharge electrodes 1,2 and the internal corona discharge electrodes 3a, and also cause contamination and a reduced lifetime of the laser system.\nFIG. 1b shows a conventional preionization unit setup wherein only one internal corona-discharge preionization electrode 3a is employed. See U.S. Pat. No. 4,240,044 to Fahlen et al. FIG. 2 shows a perspective view of a preionization unit of either of FIGS. 1a and 1b. The preionization unit includes the internal electrode 3a and the external electrode 6. The area of most intense discharge 4 is shown at the surface of the dielectric tube 3b nearest the external electrode 6.\nAnother problem with conventional corona-like preionization units is illustrated in FIG. 3. In the preionization unit of FIG. 3, an internal preionization electrode 3a is shown surrounded by a dielectric tube 3b. An external preionization electrode 6 is shown abutting the surface of the dielectric tube 3b. The dielectric tube 3b often exhibits an unsatisfactorily non-uniform surface discharge 4a in this configuration. The non-uniform surface discharge leads to instabilities such as arcing from areas of higher charge density. The lack of uniformity of surface discharge also can cause an unstable xe2x80x9cjitterxe2x80x9d of the laser output. This jitter is a fluctuation of the interval between successive laser pulses from an evolving instability in the ignition from one laser pulse to another. This variance, or jitter, is undesirable and makes laser performance legs reproducible.\nOther problems are associated with conventional corona-like preionization units such as that illustrated in FIG. 4. Some of the UV-light emanating from the outer surface of the dielectric tube 3b unit illuminates the main discharge volume 5, as is desired. However, some of the gas volume outside of the main discharge region 5 is also illuminated by the UV-light. The UV-light is preionizing a larger gas volume than is either requited or desired.\nA disadvantage related to this is illustrated in FIG. 4, which shows that at high repetition rate operation, arcing occurs across the gas volume between the external electrode of the preionizer 6 and the grounded main electrode 2. Arcing of this kind puts constraints on the maximum achievable repetition rate. Moreover, even before the onset of visible arcing of this kind takes place, the laser pulse energy is substantially reduced by parasitic discharges in the additionally preionized gas volume. These parasitic discharges produce an instability in the laser operation.\nMoreover, as may be understood from inspection of the arrows pointing away from the tube 3b of FIG. 4, some UV-light is undesirably misdirected away from the main discharge region 5 and is absorbed by the dielectric laser chamber walls. As a result, charges build up on the walls and further inefficient arcing and parasitic discharging occurs. To address this problem, Japanese Pat. Application No. 3-9582 and U.S. Pat. No. 5,337,330 to Larson each disclose a shielding element, shown as reference numerals 6 and 36, respectively, to reduce the electric field strength between the main electrode and the dielectric pipe.\nIt is accordingly an object of the invention to design a preionization unit for a laser having a high quality gas discharge by providing an intense, yet uniform, preionization of the gas volume between the main discharge electrodes.\nIt is also an object of the invention to provide a dielectric tube which prevents over-flashing and arcing at the tube ends.\nIt is another object of the invention to prevent electrical oscillations from arising out of residual energies stored in the dielectric tube,\nIt is an object of the invention to provide an external electrode which shields the walls of the discharge chamber and the gas volume outside of the main discharge area from the effects of the preionization unit.\nThe present invention meets these objects and addresses the shortcomings of conventional preionization techniques by providing a preionization device for a gas laser which comprises an internal preionization electrode having a dielectric tube around it and an external preionization electrode displaced from the dielectric housing by a small gap. The dielectric tube includes two cylindrical regions of differing outer radii of curvature. One end of the tube is open to allow an electrical connection to the internal electrode, and the other end is closed. The open end of the tube has a larger radius of curvature than the closed end. The internal electrode connects to circuitry external to the discharge chamber at the open end of the tube via a conductive feedthrough which penetrates through the housing. The external circuitry prevents voltage oscillations caused by residual energy stored as capacitance in the dielectric housing. The external preionization electrode, which is connected electrically to one of the main discharge electrodes, is formed to shield the internal preionization electrode from the other main discharge electrode to prevent arcing therebetween. The external electrode is also formed to shield the outer gas volume and walls of the discharge chamber from the preionization unit. A semi-transparent external electrode prevents electrical field distortion near the main gas discharge."}
{"text": "The neurotransmitter serotonin [5-hydroxytryptamine (5-HT)] is involved in multiple central nervous facets of mood control and in regulating sleep, anxiety, alcoholism, drug abuse, food intake, and sexual behavior. In peripheral tissues, serotonin is reportedly implicated in the regulation of vascular tone, gut motility, primary hemostasis, and cell-mediated immune responses. Walther, D. J., et al., Science 299:76 (2003).\nThe enzyme tryptophan hydroxylase (TPH) catalyzes the rate limiting step of the biosynthesis of serotonin. Two isoforms of TPH have been reported: TPH1, which is expressed in the periphery, primarily in the gastrointestinal (GI) tract; and TPH2, which is expressed in the brain. Id. The isoform TPH1 is encoded by the tph1 gene; TPH2 is encoded by the tph2 gene. Id.\nMice genetically deficient for the tph1 gene (“knockout mice”) have been reported. In one case, the mice reportedly expressed normal amounts of serotonin in classical serotonergic brain regions, but largely lacked serotonin in the periphery. Id. In another, the knockout mice exhibited abnormal cardiac activity, which was attributed to a lack of peripheral serotonin. Côté, F., et al., PNAS 100(23):13525-13530 (2003).\nBecause serotonin is involved in so many biochemical processes, drugs that affect serotonin levels are often attended by adverse effects. Thus, a need exists for new ways of treating diseases and disorders that are affected by serotonin."}
{"text": "1. Field of the Invention\nThe present invention relates to a laser driving device (laser driving circuit) and an optical apparatus such as an optical disc apparatus using such a laser driving device (a recording apparatus using an optical disc).\n2. Description of the Related Art\nA recording/reproducing apparatus using a laser as an optical source has been used in various fields. For example, optical disc recording/reproducing apparatus (hereinafter, only referred to as an optical disc apparatus) has been attracting attention.\nAs a laser used as a light source, a semiconductor laser using a semiconductor material has been widely used for various kinds of devices in recent years. This is because it is very small and responses to driving current at high speed.\nAs a writable optical disc to be used as a medium for recording or reproducing data, a phase-change optical disc, a magnetic optical disc, and so on have been widely known in the art. These optical discs perform recording, reproducing, and erasing operations by changing the strength of irradiating laser beams. Usually, when recording information on an optical disc, a so-called optical strength modulation system that forms marks and spaces by a change in intensity of a laser beam. In the recording operation, the optical disc may receive the irradiation of a laser beam at a high peak intensity as much as 30 mW or more. In the reproducing operation, the optical disc may receive the irradiation of a laser beam at an intensity lower than that of the recording operation (e.g., 1 mW) to read the information without destroying recording marks.\nMark edge recording, which uses changes at both ends of a recording mark, has become mainstream because of superiority thereof for achieving higher density of a recent writable optical disk. In the mark edge recording, however, a data error may occur by distortion in the shape of the mark. As a technique for preventing data errors due to the mark shape distortion in the mark edge recording, a write strategy technique in which recording power is modulated into multi-leveled power by pulse division (see, for example, Japanese Published Patent Application No. 2007-141406 and “Industry-leading levels of low noise and high response: Clear a wall of achieving a technology of Blu-ray with 8× speed recording/reproduction”, CX-PAL Semiconductor & Component News No. 74 (October, 2007) published by Sony Corp. (URL:http://www.sony.co.jp/Products/SC-HP/cx pal/vol74/pdf/featuring2 bd.pdf))."}
{"text": "1. Field of the Invention\nThe invention relates in general to circuit breakers and more particularly to molded case circuit breakers having an adjustable rating.\n2. Description of the Prior Art\nCircuit breakers are widely used in commerce and industry to protect electrical apparatus from over-current conditions. The varying requirements of different applications require circuit breakers in a variety of ratings. Manufacturers and suppliers are thus required to produce and hold in inventory a large number of different types of circuit breakers. Considerable savings would result if a single circuit breaker could provide a variety of ratings.\nCircuit breakers having adjustable ratings are described in U.S. Pat. No. 3,575,679 issued to James P. Ellsworth and John Zipay, assigned to the assignee of the present invention, and U.S. Pat. No. 3,593,234 issued to Allan P. Charbonneau. Circuit breakers described in the above-mentioned U.S. patents employ adjustment means having a variety of locking devices and biasing means to permit adjustment of the circuit breaker rating.\nSafety problems can result however from a circuit breaker capable of covering a range of ratings. For instance, if a circuit breaker is properly adjusted and installed on a 200 ampere line, and conditions on the line cause repeated tripping, there is a natural human tendency to increase the tripping current level of the circuit breaker to prevent the nuisance of repeated trippings. This, of course, will result in a dangerous condition if the tripping current level of the circuit breaker is increased above 200 amperes. Thus, a compromise must be established between convenience and safety.\nOne answer in the past was to provide circuit breakers having interchangeable trip units. In order to change the rating of the circuit breaker the entire trip unit was removed in the field and a new trip unit of different rating installed. This was a complex procedure, however, and sometimes resulted in improper installation of the new trip unit. This in turn resulted in failure of the circuit breaker, causing extensive damage to the circuit breaker and to the apparatus being protected. Also, the problems associated with manufacturing and stocking a variety of unit ratings remained.\nIt would be desirable to provide a circuit breaker adjustable over a range to provide a variety of ratings which would avoid safety problems associated with adjustable circuit breakers, be of simple construction, and provide for convenient adjustment."}
{"text": "The prostaglandins are a group of hormone-like substances which may be viewed as derivatives of prostanoic acid. Several prostaglandins are found widely distributed in mammalian tissue and have been isolated from this source. These prostaglandins have been shown to possess a variety of biological properties such as broncho-dilation and the ability to reduce gastric secretion.\nThe present invention concerns PGE.sub.1 and PGE.sub.2 derivatives in which the 11-position (using the prostanoic acid numbering system) is a methylene group, i.e. the 11-hydroxyl group normally present in PGE.sub.1 and PGE.sub.2 has been removed and is replaced with hydrogen. The preparation of the parent molecules of this series, 11-deoxy-PGE.sub.1 and 11-deoxy-PGE.sub.2, is reported in J. Org. Chem. 38, 951 (1973)."}
{"text": "1. Field of the Invention\nThis invention relates to a fuel cell using methanol as a fuel, and more particularly to an improvement of a fuel cell using an electrolyte ion exchange membrane.\n2. Description of the Prior Art\nA fuel cell using an electrolyte ion exchange membrane is disclosed, for example, in Japanese patent application Kokai (Laid-open) No. 56-138874, where catalyst particles are fixed to the surface of an ion exchange membrane, and the ion exchange membrane and an electrode are integrated into one body. The procedures for fixing catalyst particles to the surface of an ion exchange membrane are described in detail therein together with the prior art ones. According to one of the procedures disclosed therein, the surface of an ion exchange membrane is roughened and catalyst particles are fixed to the roughened surface, where a binder such as polytetrafluoroethylene, etc. is required when the amount of catalyst particles to be fixed exceeds about 2 mg/cm.sup.2.\nAccording to another procedure disclosed in U.S. Pat. No. 3,134,697, an electrode material comprising metal or alloy powder and a resin binder is distributed onto the surface of an ion exchange resin membrane and integrated with the ion exchange resin membrane at an elevated temperature under pressure.\nAll of the integrated electrodes disclosed in Japanese patent application Kokai (Laid-open) No. 56-138874 are hydrophilic on the electrode surfaces, and even if polytetrafluoroethylene having a water-repellent property is used as a binder, the electrode surfaces are substantially hydrophilic and not water-repellent.\nThus, the integrated electrode having catalyst particles integrated with an ion exchange membrane so far proposed is hydrophilic, and is not suitable for an oxidizer electrode (or cathode or air electrode) of a fuel cell using methanol as fuel, because, in the case of using such an oxidizer electrode, a three-phase boundary is hardly obtainable among catalyst particles (solid), a fuel (liquid) and an oxidizer (gas) with the result of a poor cell performance."}
{"text": "1. Field of the Invention\nThe invention relates to a novel washing agent and a washing method utilized for laundry, dry-cleaning, and washing of electronic parts, circuit substrates, precision machines, etc.\n2. Description of the Prior Art\nIn conventional washing methods, water has been used to remove water-soluble soil effectively in laundry since washing was begun historically. However, garments made of material apt to sustain damages such as shrinkage and loosing their trim appearance resulting from the use of water, such as wool and silk are treated with dry-cleaning. Various washing agents are properly used taking damages given to garments into consideration, for example, perchloroethylene, organic solvents etc. which are effective in removing oily soil are employed for dry-cleaning, and CFC-113 which can retain good feeling or garments, and petroleum solvents are also employed properly. Powder cleaning using powdered corn core, etc. and shampoo cleaning utilizing foams are also known."}
{"text": "1. Field of the Invention\nThis invention relates to a roasting device and, more particularly, to a coffee bean roasting device for roasting small quantities of coffee beans in the convenience of a person's home.\n2. Description of the Prior Art\nCoffee beans develop color, taste, and smell during roasting. Roasting processes must produce consistent bean roasts for a wide variety of desired bean flavors. If the beans are not circulated evenly inside the roasting chamber, coffee beans will not roast evenly. Thus, even a skilled roast master has difficulty reproducing consistent coffee roasts from time to time.\nThere are two primary methods for roasting coffee beans. The first method is drum roasting. Drum roasting involves a technique similar to clothes drying. Unroasted or green coffee beans are placed in a round drum. Heating elements adjacent to the drum heat both the drum and the air inside of the drum. The heat from the inside wall of the drum is transferred to the coffee beans. As the coffee beans roll inside the drum, the coffee beans are heated and scorched, changing from a green state to a roasted condition. To generate darker roasts, the coffee beans are heated in the drum at a higher temperature or at the same temperature but for a longer time. U.S. Pat. No. 4,860,461 to Tamaki, et al., shows a drum roaster. In Tamaki, stirring and exposing the beans to radiant heat from infrared heaters adjacent to the drum roasts the beans in the drum.\nThe second method involves a fluid bed roaster that blows air up from the bottom end of a roasting chamber. The air raises the green coffee beans up on a bed of hot air causing the beans to circulate inside the chamber. U.S. Pat. No. 3,964,175 to Sivetz shows such a fluid bed coffee roasting system. In Sivetz, air is heated to a selected temperature that roasts the beans primarily by convection. The coffee beans are secondarily roasted by the conduction of heat from bean to bean and from the inside walls of the roasting chamber to the coffee beans.\nCoffee bean roasting systems, such as those described in Tamaki and Sivetz, are very large complex machines. Thus, current roasting systems are usually operated in industrial locations. While large roasting machines may be suitable for the coffee roasting industry or for retail coffee shops that require roasting high volumes of coffee beans daily, they are ill-suited for home use that require smaller quantities of coffee beans. Moreover, few coffee lovers would tolerate a machine the size of a refrigerator taking up valuable kitchen space no matter how much they enjoy roasting their own coffee.\nExisting roasting machines are also difficult to operate and expensive to maintain. Additionally, large roasting machines require custom equipment and include complex exhaust systems. Thus, coffee roasting systems are not feasible for ordinary home use. Finally, large coffee roasting machines are expensive.\nAccordingly, a need remains for a home use coffee roasting device that evenly roasts a small quantity of coffee beans. A need also remains for a home use coffee roasting device that is simple, small, inexpensive, and easy to operate and maintain."}
{"text": "1. Field of the Invention\nThis invention regards a method of dynamic positioning of a vessel. More particularly, it regards a method that involves two or more floating buoys being coupled to a vessel by means of buoyancy, wherein the buoys are provided with positioning machinery. The buoys form part of a dynamic positioning system and are preferably provided with energy from a generator placed on the vessel.\n2. Description of Related Art\nWhen operating in relatively shallow waters offshore it is customary to use anchor based mooring methods. Anchor systems of this type reveal substantial flaws when it comes to anchoring a vessel at greater depths. The main reason for this is that the length of chain or steel cable required causes a weight-induced strain on the mooring that takes up a considerable share of the tensile strength of the mooring.\nThus it has become more and more common to equip vessels operating in deep waters with so-called dynamic positioning. Vessels such as drilling vessels, supply vessels, accommodation vessels or cargo vessels are to an increasing extent being equipped with this kind of equipment.\nIn addition to navigation and control systems, a dynamic positioning system also includes positioning machinery. It is essential that the vessel be equipped with positioning machinery that is arranged to exert a pushing force in at least two directions and in two positions on the vessel. Positioning machinery of this type typically includes several positioning thrusters that are generally arranged for rotation in the horizontal plane, so as to allow them to exert a pushing force in any horizontal direction.\nFor vessels originally built to use conventional anchor system, the conversion required to allow them to use dynamic positioning has proven to be a relatively complex and costly process."}
{"text": "The present invention is directed to cleaning agent compositions, especially useful in the electronics industry, containing acyclic aliphatic hydrofluorocarbons and a co-solvent.\nVarious organic solvents have been used as cleaning liquids for the removal of contaminants from articles and materials. Certain chlorofluorocarbons such as 1,1,2-trichloro-1,2,2-trifluoroethane have been reported as useful for this purpose, particularly with regard to cleaning organic polymers and plastics which may be sensitive to other more common and more powerful solvents such as trichloroethylene or perchloroethylene.\nSolvent cleaning and degreasing with chlorofluorocarbon-based solvents have been used extensively in industry for cleaning solid surfaces, particularly the surfaces of intricate parts. Common uses include the removal of fluxes from circuit boards and removing oils, grease and abrasives from machine parts. Cleaning procedures include exposing the cool article to vapors of the boiling chlorofluorocarbon cleaning composition and allowing any chlorofluorocarbon to evaporate from the cleaned article. Another cleaning technique involves immersing the article to be cleaned in a sump of boiling chlorofluorocarbon cleaning composition which removes most of the soil from the article, followed by a second immersion of the article in fresh chlorofluorocarbon at about room temperature and finally exposing the cooled article to vapors over the boiling sump. However, hydrofluorocarbons used alone are not effective cleaning agents.\nRecently, the long term environmental effects of chlorofluorocarbons have come under scientific scrutiny because it has been postulated that these materials because of their high stability are able to reach the stratosphere where under the influence of ultraviolet radiation release chlorine atoms which, in turn, undergo chemical reaction with stratospheric ozone. Reduction of the stratospheric ozone layer would increase the amount of ultraviolet radiation reaching the earth's surfaces. In view of the potential environmental problems associated with stratospheric ozone depletion, there is a need for new materials possessing properties which make them useful substitutes for applications in which chlorofluorocarbons have been used.\nAccordingly, there is a need in the industry for effective cleaning agents for removing oils, grease and abrasives, fluxes, water, etc., which are nonflammable, possess low toxicity and are degradable in the lower atmosphere."}
{"text": "1. Field of the Invention\nThe present invention relates to an opto-electric device.\nThe present invention further relates to a method for manufacturing an opto-electric device.\nThe present invention further relates to a semi-finished product that can be used in the method.\nThe present invention still further relates to a method of manufacturing the semi-finished product.\n2. Related Art\nAn optoelectric device is a device that provides for an optical effect in response to an electric signal, or that generates an electric signal in response to an optical stimulus. Examples of the first are light emitting diodes, such as organic light emitting diodes and electro chromic devices. Examples of the second are photo voltaic cells.\nFor large area OLED lighting on flexible plastic substrates, a large current is required to drive the system. The present thin film materials used for the anode (e.g. ITO) and cathode (e.g. Ba/Al) have a large resistivity and the large currents give rise to substantial voltage drop, which determine a non-uniform light emission. For producing large area flexible OLED devices on plastic substrates there is a need for an additional metallization structure of the plastic substrate. For reducing the manufacturing costs, such structured metallization coatings will preferably be applied on rolls of plastic foil using an inline roll-to-roll web coating process.\nAccordingly, for optoelectric devices, such as light emitting devices and electro-chromic devices, but also for photo-voltaic products there is a need for a metallization structure that on the one hand has a good electrical conductivity, while on the other hand has a high transmission for photon radiation.\nWO2007/036850 describes an organic diode device that comprises an organic diode structure having an anode layer, a cathode layer and an organic layer. One of the anode layer and the cathode layer has a set of contact areas that are distributed over a face of said structure. A barrier layer hermetically covers said structure and is provided with a set of openings aligned with said set of contact areas. A metal conductor has been electroplated on said barrier layer and contacts the set of contact areas via the set of openings.\nThe electroplated metal conductor shunts the anode and the cathode, and therewith provides for an even voltage distribution over the area of a large organic diode device and therewith an even luminance."}
{"text": "1. Field of the Invention\nAn adjustable vehicle steering column is locked to the vehicle chassis by a remote-controlled electrically-operable locking arrangement that is carried solely by a locking arm that is pivotally connected between the steering column and the vehicle chassis.\n2. Description of the Related Art\nIt is well known in the prior art to provide vehicle steering columns wherein the position of the steering wheel can be adjusted in terms of height and/or in terms of depth. The position of the steering wheel is selected by the driver and is maintained in place by a locking system whose manual control consists of an operating lever that is worked by the driver.\nThese different types of locking devices, operated by means of a lever, possess several inconveniences. First of all, the lever can be positioned only in the proximity of the locking axis, in other words, in an area where accessibility is reduced along with poor visibility, which creates a risk of injury as the fingers are squeezed. Moreover, the locking and unlocking operations necessitate a major effort and is interfered with by the driver's knees. Finally, the operating lever is situated in areas that must be used for other functions so as to provide a good ergonomy for the driver's position. This environment can be the area of computers when the locking axis is placed adjacent or above the column spindle. This environment can also be the area where the driver's legs are located when the locking axis is placed below the column shaft, which introduces a risk of impact against the knees in case of an accident.\nThe present invention was developed to provide a device for locking an adjustable automotive vehicle steering column, which will eliminate the above-described inconveniences by providing remote locking control. This remote control must employ a simple and rather small device to allow great adaptation flexibility depending on the type of steering column and the environment of the driver's position."}
{"text": "This invention relates to an infrared gas burner plate used for a radiation type gas burner.\nAs for a burner plate of this kind, there has been hitherto known such a type as that shown in FIG. 1, for instance. A sintered plate body a composed chiefly of a heat resisting material such as ceramic or the like is provided at its front surface with a large number of grooves b crossing one another and a large number of convex portions c defined by these grooves b The plate body a is additionally provided with a large number of burner holes d piercing therethrough between the front and rear surfaces thereof and distributed over the grooves b and the convex portions c so that each of convex portions c may be made red-hot by being heated through the burner holes d made therein and the burner holes d in the surrounding grooves b around the same, and thereby the burner plate may be improved in radiation efficiency. This conventional type of burner plate, however, is such that the same is liable to be broken along any of the grooves b when subjected to a bending force. Moreover, it takes only about 40 seconds from ignition until the ignition can be confirmed visually from a red-hot condition of a corner portion of any of the convex portions c that is raised in temperature most easily. This makes it easy to more rapidly confirm the ignition as compared with a conventional flat plane type plate in which it takes one minute to confirm the ignition. However, it is desirable that this time be shorter so that use of a burner plate of this kind may be further facilitated.\nThis invention has for its object to provide a burner plate which has an improved bending strength, and in which the time necessary for confirmation of the ignition is shortened, and also has an improved radiation efficiency. More particularly, it is an object to provide a burner plate of the type wherein a sintered plate body composed chiefly of a heat resisting material such as ceramic or the like is provided at its front surface with a large number of grooves crossing one another and a large number of convex portions defined by these grooves, the plate body being additionally provided with a large number of burner holes piercing therethrough between the front and rear surfaces thereof and distributed over the grooves and the convex portions characterized in that each of the intersecting portions of these grooves is provided with a small convex portion having at its center the burner hole. In a preferred embodiment, each of the small convex portions is frusto-conical in form."}
{"text": "The invention herein relates to an automobile steering column switch comprising a support, which can be mounted to a steering column tube of an automobile steering column, said support having arranged on it at least one wiper switch and/or blinker switch and/or driving light switch.\nAs has been known, motor vehicles consist of a plurality of complex construction assemblies, whereby the steering column switch assembly to be installed on the steering column tube of the motor vehicle represents an important feature.\nTypically, steering column switches are mounted in radial direction relative to the steering column on the steering column\"\"s steering column tube, thereby rendering the steering column switch immobile relative to the steering wheel when the steering wheel and hence the steering column guided inside the steering column tube are moved. As a rule, the steering column switch comprises a blinker switch and a driving light switch, in which case the driving light switch frequently is configured as a rotary switch being an integral part of the switch lever of the blinker switch. Furthermore, the steering column switch frequently comprises a wiper switch for graduated activation of the windshield wipers and activation of the switch functions for the windshield washer system. For example, such steering column switches have been known from DE Auslegeschrift 2,810,790 and DE Patent 3,532,532 C2.\nEvery steering column switch comprises a plurality of interacting components that perform a large number of switch functions. These components are accommodated in an extremely small space and, therefore, make the construction of the steering column switch complex. Therefore, depending on the type of vehicle, individual components are more or less pre-assembled and subsequently mounted to a support located on the steering column tube of the steering column.\nA similar situation exists in the case of an ignition starter switch with integrated steering lock because the ignition starter switch frequently performs two additional switch functions such as, for example, activation of an immobilizer cooperating with a transponder integrated in the ignition key. As a result, ignition starter switches are becoming more and more complex because integrated code-recognition systems are to protect the vehicle against unauthorized use.\nThe two assemblies, namely the steering column switch and the ignition starter switch, are mounted separately to the steering column tube at the time the motor vehicle is assembled, whereby, in most cases, the ignition starter switch is located at a distance from the steering column switch. However, considering the viewpoint of technology, the automobile industry has demanded reduced assembly times for motor vehicles, in particular on the assembly line.\nThe problem to be solved by the invention herein is to provide an automobile steering column switch of the above-described type, in which case said switch can be pre-assembled practically completely and can incorporate the switch functions of an ignition starter switch.\nIn accordance with the invention herein, this problem has been solved in that\nthe support comprises an integrated printed circuit board provided with electrical contacts to which the switch functions relating to the wiper switch and/or blinker switch and/or driving light switch and an ignition starter switch can be transferred,\nsaid support is configured in one piece with an accommodation space for the ignition starter switch, and\nsaid support can be mounted, together with the switches, as pre-assembled module to the steering column tube.\nConsequently, the support of the steering column switch not only accommodates the printed circuit board for performing the switch functions of the steering column switch and the ignition starter switch but also is configured in such a manner that it can accommodate the ignition starter switch. Therefore, all of the switch functions combined are transferred as one single function to the printed circuit board provided in the steering column switch. As a result, it is no longer necessary to install a plurality of electrical cables in the steering column switch, said cables taking up an unnecessary amount of space and being prone to trouble. In addition, the electrical cable connections of the ignition starter switch are eliminated because its switch functions are now transferred directly to the printed circuit board on the support. The particular advantage of this is that the pre-assembled ignition starter switch can be inserted in one piece in the respective accommodation space on the support of the already pre-assembled steering column switch. The steering column switch may then be mounted in the form of a module to the steering column of the motor vehicle. Consequently, this reduces the number of required components, the required assembly facilities, as well as the assembly time on the assembly line during the preliminary and final assemblies of the motor vehicle. Therefore, a separation of mechanical and electrical assembly functions has been created.\nAs an alternative, this problem has also been solved in thatxe2x80x94the support comprises a printed circuit board with electrical contacts, to which the switch functions relating to the wiper switch and/or the blinker switch and/or the driving light switch, as well as an ignition starter switch, can be transferred,\nthe support can be mounted as pre-assembled module to the steering column tube, and\nthe ignition starter switch can be mounted as a pre-assembled module to the steering column tube, separately from the support, by means of a clamp connection.\nConsidering the printed circuit board and the combination of the switch functions of the steering column switch in conjunction with the ignition starter switch, this alternative solution offers the same advantages as the first solution. While in this case the ignition starter switch is not directly integrated in the support of the steering column switch, but rather mounted as a separately pre-assembled component on the steering column tube of the steering column in close proximity of the steering column switch, the switch functions of the ignition starter switch can be transferred directly to the printed circuit board mounted to the support due to its close proximity to this support. In order to simplify the assembly effort relating to the ignition starter switch provided as separate module and to allow fast assembly, the clamp connection consists of a clamp overlapping the steering column tube and a corresponding counter-clamp adjoining the ignition starter switch. In this case the clamp overlapping the steering column tube is connected with the counter-clamp on the ignition starter switch by means of screws. Due to this, only a few hand movements are required to mount the ignition starter switch quickly to the steering column tube.\nIn the case of an advantageous embodiment of the invention herein the wiper switch and/or blinker switch and or driving light switch are mounted to the support by means of clip connections. These become a fixed integral part of the support simply by inserting the switch in the respective accommodation space on the support. As a result, these clip connections may be released conveniently for repair and maintenance purposes, thereby permitting the quick and easy removal of the switches from their respective accommodation spaces.\nIn order to reduce assembly efforts even further, one modification of the invention herein provides that the support be mounted in position to the steering column tube by means of a clip connection or bayonet connection, whereby the support rests on a ball-bearing located on the steering column. This ball-bearing is placed in the support of the steering column switch during preliminary assembly and placed on the steering column at the time of assembly in such a manner that, due to the clip connection or bayonet connection between the support and the steering column tube, the steering column switch is mounted rigidly on said steering column tube.\nIn accordance with another modification of the solutions offered by the invention herein the electrical contacts of the printed circuit board are configured as micro-switches. These micro-switches can have multiple functions, be subjected to frequent use, and are resistant to mechanical stress and distinguished by high reliability. As an alternative, as has been well known, the said electrical contacts of the printed circuit board may be configured as switch contacts or the like.\nIn order to be able to convert switch functions transferred to the printed circuit board in an effective manner, the circuit board is connected with a computer-controlled electronic control system of the automobile. In addition, the printed circuit board is covered with a lid that is mounted to the support. As a result of this, the printed circuit board is protected against outside influences such as damage caused by moisture and dust.\nFrequently, switch contacts located inside the steering column switch are subject to corrosion or breakage due to excessive use. To prevent this and achieve greater reliability, the ignition starter switch preferably comprises a cam control activated by means of a key, said cam control acting on the electrical contacts of the printed circuit board. Therefore, different configurations of the cams activate different switch functions that are transferred to the printed circuit board in the support of the steering column switch.\nIt is useful if the support incorporates an electrical transfer system for an airbag device. As a rule, this type of transfer system comprises a helical volute buffer spring used as electrical connection element between two parts, which move relative to each other; this is the case with a steering wheel containing an airbag and moving relative to the steering column switch. In as much as the electrical transfer system is also an integral part of the pre-assembled steering column switch module, the assembly effort is reduced even further."}
{"text": "1. Field of the Invention\nThe inventions disclosed and taught herein relate generally to induced or forced air filtration; and more specifically related to a system for removing contaminants from air used to ventilate an equipment room on an offshore drilling rig.\n2. Description of the Related Art\nIn the drilling industry, drilling fluids, such as drilling mud, which may be a mixture of clay, water, and various additives, are pumped down through the drill string into the well. The drilling fluid exits the drill string at the bottom through holes or jets in a drill bit. The fluid picks up cuttings and other solids from the well and carries them away from the bit and out of the well in the annulus between the well and the drill string.\nOnce the drilling fluid is returned to the surface, equipment, such as mud pits and shakers, are used to remove large amounts of coarse, drilling cuttings from the drilling fluids so that the drilling fluids may be reused or recirculated within the well. The cuttings are generally filtered out of the drilling fluid on a series of filters or screens provided within the shaker. The drilling fluid that has passed through the equipment is now clean of coarse debris and is ready to be cleaned via other processes prior to being ready for reuse.\nOn offshore drilling rigs, where area or space is a precious commodity, the equipment is typically located below the main work deck in an enclosed area. Drilling fluids typically contain a measure of volatile organic compounds and fugitive contaminants from the drilling process. Because of this, the enclosed room is subject to forced ventilation to remove unpleasant and/or hazardous fumes. Depending on where the fumes are vented on the rig and the various and changing environmental conditions, the fumes may infiltrate other workspaces on the rig.\nThe inventions disclosed and taught herein are directed to a filter system for removing or reducing equipment room contamination of ventilation air."}
{"text": "Convolutional neural networks (CNNs) may be used in machine learning implementations in applications such as image and video recognition, or natural language processing. A CNN may typically comprise a number of layers of computational constructs called “neurons.” Each neuron may receive as input a portion of the receptive field (e.g. the output of the previous layer), and output a value for a selected feature. A CNN may be trained based on a set of training data which provides inputs and correct outputs."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for regulating the exposure time of a light sensor. More particularly, the invention relates to a process for regulating the exposure time of a light sensor according to the variation of luminosity in the work environment in which the sensor is located.\n2. Discussion of the Prior Art\nThroughout the present description and the attached claims, with the term: “light sensor”, it is intended to indicate a light-sensitive optical element capable of converting a light signal, contained in a incident light beam, in an analogic electrical signal proportional to the light intensity of the light itself. In particular, the light sensor transforms the light image of an illuminated object into an electronic image thereof, in order to allow the carrying out of further operations of image processing, such as for example those connected to the reading of an optical code located on the object.\nIn the following, reference to light sensors of known type will be made, for example CCD or CMOS sensors, both linear and of matrix type.\nAs is known, a drawback associated with light sensors at present available on the market is correlated to the fact that their operation is strongly influenced by the variations in luminosity of the surface on which the image to be acquired is to be found (such surface can be more or less clear and more or less illuminated). In particular, the resolution of the image acquired by the sensor (that is the richness of details detected in the image acquired by it) considerably varies according to the variation of the luminosity of the surrounding environment and of the surface to be acquired: highly illuminated environments can as such create conditions of strong overexposure or saturation (little contrast) of the image on the sensor, whilst badly illuminated environments can create conditions of strong under-exposure (excessive dark) in the image on the sensor. In both cases, the image acquired by the sensor appears badly exposed and not very clear, therefore making it impossible to detect details having small sizes.\nIn particular, with reference to the reading of an optical code, the acquisition on the sensor of an image of an optical code having a low resolution is often undesired because it renders the subsequent operations of localization, reading and decoding of the code itself very approximate and rough; this is absolutely unacceptable where it is necessary to carry out a correct decoding of the code in order to unequivocally identify the object having the code thereon.\nThe drawbacks associated with the acquisition of an image having low resolution can be overcome by suitably regulating the exposure time of the sensor according to the variations of luminosity of the surrounding work environment and of the surface to be acquired. In particular, where an increase in luminosity occurs in the work environment, it would be advisable to set a shorter exposure time on the sensor, in order to avoid the risk of generating conditions of saturation of the acquired image; viceversa, in the case where a decrease in luminosity occurs in the work environment, it would be advisable to set a higher exposure time on the sensor, in order to avoid the risk of generating conditions of under-exposure of the acquired image.\nThe majority of light sensors at present available on the market allow to set an exposure time by choosing it between a range of values, more or less extensive, each differing the one from the other by descrete amounts; such values are stated by the manufacturers of the sensors."}
{"text": "This invention relates to a magnetic apparatus for producing movement, preferably linear movement which is produced by intermittent actuation of electromagnets positioned in an aligned set of relatively movable paramagnetic members.\nThe prior art is replete with disclosures of linear actuators wherein electromagnets are periodically energized to provide linear movement. Typical examples of the prior art are discussed in and cited against the earlier patent applications mentioned above.\nA particular objective of the present invention is to provide a system which is superior to prior art devices in that it will provide a high efficiency and a low weight-to-power ratio.\nThe invention disclosed herein is an improvement over my earlier devices in two primary respects. First, rather than providing electrical windings on all of the paramagnetic members in the system, the present system utilizes axially movable reactor members which are positioned between the electromagnets and do not themselves create a magnetic field. Second, the invention calls for the reduction of electrical consumption by providing for the gradual reduction of the magnetomotive force of the electromagnets while the gap is being closed by the action of such electromagnets. This reduction of magnetomotive force, while conserving energy, does not significantly reduce the force output of the device since, as the gaps between the members are closing, a lesser magnetomotive force is required to produce any given force of magnetic attraction.\nMore specifically, the invention relates in one sense to an apparatus which has a pair of relatively movable paramagnetic members which are separated by a gap which opens and closes as the members move relatively away from and toward each other. One of the paramagnetic members is provided with electrical windings which, when energized, produce a magnetic field which will pull the members together and close the gaps between them. The windings are energized by means which causes an electrical current to flow through them to create a magnetomotive force, and control means are provided for gradually reducing the magnetomotive force while the gap is being closed by the field. Preferably, the windings associated with a particular paramagnetic member include a plurality of winding sections, and the control means is constructed and operated to energize simultaneously and deenergize sequentially the sections of windings. The control means preferably is a timer provided with a set of cams which rotate together at a constant angular velocity. A set of electrical switches are sequentially operated by the cams and are connected to respective sections of windings. A single rotary cam may be operatively engaged with switches associated with the windings of two different paramagnetic members.\nIn another sense, the invention pertains to a system which includes an axially aligned set of paramagnetic members which include a central paramagnetic member and a spaced pair of paramagnetic members which are spaced apart by a constant distance and are located on opposite sides of the central member. Two intermediate paramagnetic members are located where each lies between the central paramagnetic member and one of the spaced pair of members. The energizing means is operable in a first mode to produce movement in a first direction by creating magnetic attraction between one intermediate member, the central member and one of the spaced pair of members. The energizing means is operable in a second mode to produce movement in an opposite direction by creating magnetic attraction between the other intermediate member, the central member and the other one of the spaced pair of paramagnetic members.\nIn another sense, the invention relates to a set of aligned paramagnetic members which includes a driven paramagnetic member located between two stationary paramagnetic members. A first movable intermediate paramagnetic member is located between one stationary member and the driven member; and, a second movable intermediate paramagnetic member is located between the other stationary member and the driven member. The intermediate members are on opposite sides of the driven member. The energizing means associated with this apparatus is operable in a first mode to create magnetic attraction between the first intermediate member, the driven member and the first stationary member to move the driven member in a first direction. The energizing means is operable in a second mode to create magnetic attraction between the second intermediate member, the driven member and the second stationary member to move the driven member in a second direction which is opposite to the first direction. Means are provided for repeatedly shifting the energizing means between the first mode and the second mode to move the driven member alternately in the first and second directions.\nAlthough the invention may take many forms, the presently preferred embodiments thereof are described in the following text and illustrated in the drawings."}
{"text": "Color transfer thermal printers use a color donor member which may be a sheet, but usually is in the form of a web advanced from a supply roll to a take-up roll. The color donor member passes between a printhead and a dye receiver member. The thermal printhead comprises a linear array of resistive heat elements. In operation, the resistive heat elements of the printhead are selectively energized in accordance with data from a printhead control circuit. As a result, the image defined by the data from the printhead control circuit is placed on the receiver member.\nA significant problem in this technology is that the color donor members used to make the thermal prints are generally intended for single (one time) use. Thus, although the member has at least three times the area of the final print and contains enough colorant to make a solid black image, only a small fraction of the color is ever used.\nAfter printing an image, the color donor cannot be easily reused, although this has been the subject of several patents. The primary reason that inhibits reuse of the color donor is that the color transfer process is very sensitive to the concentration of the colorant in the donor layer. During the first printing operation, color is selectively removed from the layer thus altering its concentration. In subsequent printings, regions of the donor which had been previously imaged have a lower transfer efficiency than regions which were not imaged. This results in a ghost image appearing in subsequent prints.\nThe cost associated with having a single use donor ribbon is large because of the large area of ribbon required, as well as the large excess of colorant coated on the donor member. While this technology is able to produce high quality continuous tone prints, it is desired to provide an approach which has all of the good attributes of thermal color transfer imaging but without the limitations associated with single use donor members.\nSome work has been done by others to accomplish similar goals. For example, U.S. Pat. No. 5,286,521 discusses a reusable wax transfer ink donor ribbon. This process is intended to provide a dye donor ribbon that may be used to print more than one page before the ribbon is completely consumed. U.S. Pat. No. 4,661,393 describes a reusable ink ribbon, again for wax transfer printing. U.S. Pat. No. 5,137,382 discloses a printer device capable of re-inking a thermal transfer ribbon. However, again the technology is wax transfer rather than dye transfer. In the device, solid wax is melted and transferred using a roller onto the reusable transfer ribbon.\nU.S. Pat. No. 5,334,574 describes a reusable dye donor ribbon for thermal dye transfer printing. This reusable ribbon has multiple layers containing dye which limit the diffusion of dye out of the donor sheet. This enables the ribbon to be used to make multiple prints. In addition, the ribbon may be run at a slower speed than the dye receiver sheet, enabling additional utilization. U.S. Pat. No. 5,118,657 describes a multiple use thermal dye transfer ink ribbon. This ribbon has a high concentration dye layer on the bottom and low concentration dye layer on the top. The low concentration dye layer meters or controls dye transfer out of the ribbon. This enables the ribbon to be used multiple times. U.S. Pat. No. 5,043,318 is another example of a thermal dye transfer ribbon which can be used multiple times.\nAnother problem with the resistive head thermal printers described above is high the large amount of energy used by thermal resistive printing. This is because the heat generated by resistive printing head must first heat the support layer before the heat reaches the color donor layer. While the color donor layer can be quite thin and therefore have a low heat capacity, the support layer must be relatively thick, particularly if the support is to be used for multiple printing events without tearing or distorting from the insult of the printing operation.\nYet another problem of the printers described above is that of low resolution images. As the heat from the printhead moves through the support toward the color donor layer, the heat spreads in all directions, so that the spot of color printed is necessarily larger than when it is first generated at the print head. This means the final image is of lower resolution then might be possible with a thinner color donor web."}
{"text": "Programmable logic arrays are well known. U.S. Pat. No. 4,124,899, for example, describes both the background and the uses of programmable logic arrays and field programmable logic arrays. U.S. Pat. No. 4,717,912 discloses an improved version of the programmable logic array described in U.S. Pat. No. 4,124,899 wherein an output cell is provided to allow a single circuit to be configured to simulate a programmable logic array having a combinatorial unregistered output lead or a registered output lead.\nAs is known in the programmable logic array art, one type of programmable logic array uses what is called an \"X-type\" output structure where the output signal is passed through an exclusive OR gate and stored in a register (See FIG. 2). Another type of programmable logic circuit uses what is called an \"L-type\" output structure and has a combinatorial unregistered output (See FIG. 1).\nA typical logic output cell for a programmable logic array of the \"L-type\" is shown in FIG. 1. In FIG. 1, the seven logic signals PT0-PT6 are applied to OR gate 10. The output signal from OR gate 10 optionally is passed through an exclusive OR gate (not shown in FIG. 1) which acts as either a buffer or a programmable invertor for the signal before it reaches tristate output buffer 12. The output signal from buffer 12 is sent to terminal 14. Terminal 14 can be used either as an input path or an output path depending upon the state of tristate buffer 12. If tristate buffer 12 presents a high impedance in response to logic signal PT7 being low level, then terminal 14 functions as an input terminal and the input signal on terminal 14 is transmitted on leads 19 and 17 to buffer 13. Buffer 13 provides on output leads 18b and 18a, true and complemented versions, respectively, of the signal on lead 17 and terminal 14. The structure shown in FIG. 1 provides a logical combinatorial unregistered output signal on lead 19 when buffer 12 acts as a transmission gate for passing the output signal from OR gate 10 to output terminal 14. In this context, lead 17 and buffer 13 are part of a feedback path capable of providing the true and complement of the output signal on terminal 14 to selected points within the programmable logic array.\nFIG. 2 illustrates an X-type output structure from a programmable logic array. As shown in FIG. 2, four input logic signals PT0-PT3 are applied in pairs to OR gates 31a and 31b. Logic signals PT0 and PT1 are applied to the two input leads to OR gate 31a and logic signals PT2 and PT3 are applied to the two input leads to OR gate 31b. The output leads from OR gates 31a and 31b are applied to the two input leads to exclusive OR gate 32. The output signal from exclusive OR gate 32 is applied to register 33 (typically a D-type flip-flop). The Q output signal from register 33 is passed through invertor 34 to the output terminal 36. The Q output lead from register 33 is fed back through buffer 35 which provides on output leads 37a and 37b, true and complementary versions of the signal on lead 33b, respectively.\nThe architectures of the output logic circuits with existing programmable logic arrays provide a number of different options to the designer. However, in many cases the use of one option as opposed to another requires the use of a different part or product. This can be a disadvantage if the user wishes to change the design after acquiring a substantial inventory of one type of product."}
{"text": "Many wireless communication devices include hardware modules and/or software modules to provide support for Global Navigation Satellite Systems (GNSS). The United States government provides a Global Positioning System (GPS) that typically uses 31 satellites orbiting the Earth broadcasting specialized GPS signals in a 1.6 GHz frequency band. Other global positioning systems such as the Russian GLONASS and the Chinese Beidou systems function in approximately the same manner. The term GNSS can be used to refer to satellite based positioning systems in general, and the terms GNSS and GPS can be used synonymously herein to indicate satellite based global positioning systems and signals without loss of generality.\nWireless communication device based GNSS systems can try to determine a position of the wireless communication device by receiving GNSS satellite signals. Oftentimes, the position of the wireless communication device can be presented to a user of the wireless communication device through a map displayed on the wireless communication device or on a display to which the wireless communication device is coupled. The wireless communication device can attempt to determine a position based on information obtained while in poor GNSS operating conditions, which can adversely affect the accuracy of position determination and information provided for map location display. GNSS signals are prone to fade, can suffer multipath or may be blocked by structural objects such as tunnels and/or tall buildings. In these cases, a GNSS determined position may be inaccurate.\nTherefore, what is desired is an enhancement of the performance of GNSS determined positioning, particularly when the wireless communication device operates in conditions that can result in poor reception of GNSS signals."}
{"text": "1. Field of the Invention\nThe systems and methods of this invention generally related to communications systems. In particular, the systems and methods of this invention relate to providing a variable state length initialization.\n2. Description of Related Art\nMulticarrier modulation, which is also known as Discrete Multitone Transmission (DMT), transceivers step a through a number of initialization states prior to entering steady-state communication or “showtime.” In particular, these various initialization states include channel discovery, transceiver training, channel analysis, and the like. These various initialization states allow, for example, the determination of transmitter power levels, line characteristics, training of receiver function such as equalizers or echo cancellers, or any other feature necessary to establish communication, or to exchange parameters and settings, between transceivers."}
{"text": "The farnesoid-X-receptor (FXR) is a member of the nuclear hormone receptor superfamily of transcription factors. FXR was originally identified as a receptor activated by farnesol, and subsequent studies revealed a major role of FXR as a bile acid receptor [Makishima, M., Okamoto, A. Y., Repa, J. J., Tu, H., Learned, R. M., Luk, A., Hull, M. V., Lustig, K. D., Mangelsdorf, D. J. and Shan, B. (1999) Identification of a nuclear receptor for bile acids. Science 284, 1362-5]. FXR is expressed in liver, intestine, kidney, and the adrenal gland. Four splice isoforms have been cloned in humans.\nAmong the major bile acids, chenodeoxycholic acid is the most potent FXR agonist. Binding of bile acids or synthetic ligands to FXR induces the transcriptional expression of small heterodimer partner (SHP), an atypical nuclear receptor family member that binds to several other nuclear hormone receptors, including LRH-1 and LXR alpha and blocks their transcriptional functions [Lu, T. T., Makishima, M., Repa, J. J., Schoonjans, K., Kerr, T. A., Auwerx, J. and Mangelsdorf, D. J. (2000) Molecular basis for feedback regulation of bile acid synthesis by nuclear receptors. Mol Cell 6, 507-15]. CYP7A1 and CYP8B are enzymes involved in hepatic bile acid synthesis. FXR represses their expression via activation of the SHP pathway. FXR directly induces the expression of bile acid-exporting transporters for the ABC family in hepatocytes, including the bile salt export pump (ABCB11) and the multidrug resistance associated protein 2 (ABCC2) [Kast, H. R., Goodwin, B., Tarr, P. T., Jones, S. A., Anisfeld, A. M., Stoltz, C. M., Tontonoz, P., Kliewer, S., Willson, T. M. and Edwards, P. A. (2002) Regulation of multidrug resistance-associated protein 2 (ABCC2) by the nuclear receptors pregnane X receptor, farnesoid X-activated receptor, and constitutive androstane receptor. J Biol Chem 277, 2908-15; Ananthanarayanan, M., Balasubramanian, N., Makishima, M., Mangelsdorf, D. J. and Suchy, F. J. (2001) Human bile salt export pump promoter is transactivated by the farnesoid X receptor/bile acid receptor. J Biol Chem 276, 28857-65]. FXR knockout mice have impaired resistance to bile acid-induced hepatotoxicity and synthetic FXR agonists have been shown to be hepatoprotective in animal models of cholestasis [Liu, Y., Binz, J., Numerick, M. J., Dennis, S., Luo, G., Desai, B., MacKenzie, K. I., Mansfield, T. A., Kliewer, S. A., Goodwin, B. and Jones, S. A. (2003) Hepatoprotection by the farnesoid X receptor agonist GW4064 in rat models of intra- and extrahepatic cholestasis. J Clin Invest 112, 1678-87; Sinal, C. J., Tohkin, M., Miyata, M., Ward, J. M., Lambert, G. and Gonzalez, F. J. (2000) Targeted disruption of the nuclear receptor FXR/BAR impairs bile acid and lipid homeostasis. Cell 102, 731-44]. These data show that FXR protects hepatocytes from bile acid toxicity by suppressing both cellular synthesis and import of bile acids and stimulating their biliary excretion.\nThe process of enterohepatic circulation of bile acids is also a major regulator of serum cholesterol homeostasis. After biosynthesis from cholesterol in the liver, bile acids are secreted with bile into the lumen of the small intestine to aid in the digestion and absorption of fat and fat-soluble vitamins. The ratio of different bile acids determines the hydrophilicity of the bile acid pool and its ability to solubilize cholesterol. FXR activation increases the hydrophilicity of the pool, decreasing the intestinal solubilization of cholesterol, effectively blocking its absorption. Decrease absorption would be expected to result in lowering of plasma cholesterol levels. Indeed direct inhibitors of cholesterol absorption such as ezetimibe decrease plasma cholesterol, providing some evidence to support this hypothesis. However ezetimibe has limited efficacy which appears due to feedback upregulation of cholesterol synthesis in cells attempting to compensate for depletion of cholesterol. Recent data have shown that FXR opposes this effect in part by directly repressing the expression of HMGCoA reductase via a pathway involving SHP and LRH1 [Datta, S., Wang, L., Moore, D. D. and Osborne, T. F. (2006) Regulation of 3-hydroxy-3-methylglutaryl coenzyme A reductase promoter by nuclear receptors liver receptor homologue-1 and small heterodimer partner: a mechanism for differential regulation of cholesterol synthesis and uptake. J Biol Chem 281, 807-12]. FXR also decreases hepatic synthesis of triglycerides by repressing SREBP1-c expression by an alternate pathway involving SHP and LXRalpha. Thus compounds which modulate FXR activity may show superior therapeutic efficacy on plasma cholesterol and triglyceride lowering than current therapies.\nMost patients with coronary artery disease have high plasma levels of atherogenic LDL. The HMGCoA reductase inhibitors (statins) are effective at normalizing LDL-C levels but reduce the risk for cardiovascular events such as stroke and myocardial infarction by only about 30%. Additional therapies targeting further lowering of atherogenic LDL as well as other lipid risk factors such as high plasma triglyceride levels and low HDL-C levels are needed.\nA high proportion of type 2 diabetic patients in the United States have abnormal concentrations of plasma lipoproteins. The prevalence of total cholesterol>240 mg/dl is 37% in diabetic men and 44% in diabetic women and the prevalence for LDL-C>160 mg/dl are 31% and 44%, respectively in these populations. Diabetes is a disease in which a patient's ability to control glucose levels in blood is decreased because of partial impairment in the response to insulin. Type II diabetes (T2D), also called non-insulin dependent diabetes mellitus (NIDDM), accounts for 80-90% of all diabetes cases in developed countries. In T2D, the pancreatic Islets of Langerhans produce insulin but the primary target tissues (muscle, liver and adipose tissue) develop a profound resistance to its effects. The body compensates by producing more insulin ultimately resulting in failure of pancreatic insulin-producing capacity. Thus T2D is a cardiovascular-metabolic syndrome associated with multiple co-morbidities including dyslipidemia and insulin resistance, as well as hypertension, endothelial dysfunction and inflammatory atherosclerosis.\nThe first line treatment for dyslipidemia and diabetes is a low-fat and low-glucose diet, exercise and weight loss. Compliance can be moderate and treatment of the various metabolic deficiencies that develop becomes necessary with, for example, lipid-modulating agents such as statins and fibrates, hypoglycemic drugs such as sulfonylureas and metformin, or insulin sensitizers of the thiazolidinedione (TZD) class of PPARgamma-agonists. Recent studies provide evidence that modulators of FXR may have enhanced therapeutic potential by providing superior normalization of both LDL-C and triglyceride levels, currently achieved only with combinations of existing drugs and, in addition, may avoid feedback effects on cellular cholesterol homeostasis.\nThe compounds of the present invention bind to and selectively modulate FXR very efficiently, resulting in reducing cholesterol absorption lowering LDL cholesterol and triglycerides and reducing inflammatory atherosclerosis. Since multiple facets of combined dyslipidemia and cholesterol homeostasis are addressed by FXR modulators, they are expected to have an enhanced therapeutic potential compared to the compounds already known in the art."}
{"text": "It is typical of the usual tractor-drawn plow that it is connected or supported at its laterally spaced front portions by vertically shiftable hitch means on a tractor and on vertically shiftable linkage means on a furrow wheel. U.S. Pat. No. 3,236,313 discloses an arrangement in which the furrow wheel linkage includes a two-way hydraulic motor and the linkage is mechanically connected to the hitch means so that operation of the motor activates the furrow wheel linkage and hitch means substantially in unison so that substantial leveling of the plow frame is achieved during movement between raised and lowered positions. It is also typical of such plow that, during lowering of the frame from fully raised or transport position to a position for opening up the field, the frame will descend to a greater distance in relation to the front furrow wheel than when it is lowered from transport to the normal plowing position, because, in the opening-up plowing position, the furrow wheel is traveling on unplowed land and, in normal position, the furrow wheel runs in a previously formed furrow. In the recent past it has been necessary to control a hydraulic cylinder in two phases, one stroke for lowering to open-up plowing position and a different stroke for lowering to normal plowing position."}
{"text": "Conventionally, as a method for forming quantum dots, the forming method based on Stranski-Krastanow mode (S-K mode) is known.\nS-K mode is a mode in which a semiconductor crystal to be epitaxially grown grows two-dimensionally (film growth) at the start of the growth and grows three-dimensionally on the stage where the semiconductor crystal has exceeded the elastic limit of the film. A film having a larger lattice constant than the base material is epitaxially grown, whereby quantum dots of three-dimensionally grown islands are self-assembled.\nS-K mode can easily make quantum dots self-assembled, and is widely used in the field of photosemiconductor devices, etc.\nRecently, techniques in the fields of quantum information and quantum computation attract a great deal of attention. In these fields it is very important to form a quantum dot in a prescribed position in a prescribed size.\nHowever, the conventional method for forming quantum dots described above makes positions and sizes of the formed quantum dots random.\nAs techniques of controlling positions for quantum dots to be formed in, the following techniques have been proposed.\nFor example, Japanese published unexamined patent application No. 2000-315654, the specification of U.S. Pat. No. 5,229,320 and Appl. Phys. Lett., Vol. 76, No. 2, p. 167-169, (2000) propose techniques of controlling positions for quantum dots to be formed in by forming in advance by means of electron beams concavities in the surface of a semiconductor substrate in positions for quantum dots to be formed in. When a semiconductor layer is grown on a semiconductor substrate with concavities formed in, the semiconductor layer tends to grow at a higher rate in the concavities, whereby quantum dots can be formed in the concavities.\nAppl. Phys. Lett., Vol. 75, No. 22, p. 3488-3490, (1999) and Phys. Stat. Sol. (b) 224, No. 2, p. 521-525, (2001) propose a technique of forming in advance by means of an STM (Scanning Tunneling Microscope) deposits in positions for quantum dots to be formed in to thereby control the positions for the quantum dots to be formed in. When a semiconductor layer is grown on a substrate with deposits formed on, concavities are formed in the semiconductor layer surface above the deposits. When another semiconductor layer is grown on the semiconductor layer with the concavities formed in, quantum dots are formed in the concavities because said another semiconductor layer grows in the concavities at a faster growth rate.\nAppl. Phys. Lett., Vol. 77, No. 16, p. 2607-2609, (2000) proposes a technique of forming trenches in the surface of a semiconductor substrate by means of an AFM (Atomic Force Microscope) to form quantum dots in the trenches.\nHowever, in the above-described technique of forming the concavities by means of electron beams, the surface of the semiconductor substrate is contaminated with carbon when the concavities are formed by means of electron beams, which makes it difficult to form quantum dots of good quality.\nIn the above-described technique of controlling positions of quantum dots by using an STM, crystal defects take place in the semiconductor layer with the concavities formed in. It is difficult to form quantum dots of good quality on the semiconductor layer having the crystal defects. Furthermore, the diameter of the concavities are somewhat enlarged, and it is difficult to form fine quantum dots of a below 40 nm (including 40 nm)-diameter.\nIn the technique of controlling positions of quantum dots by using an AFM, the trenches are formed in the substrate surface by mechanically scraping the substrate surface with the probe of the AFM. Therefore, the configurations of the trenches are non-uniform. It is difficult to form quantum dots in a uniform size. The trenches are formed by mechanically scraping the substrate surface, which causes crystal defects in the trenches. It is difficult to form quantum dots of good quality without crystal defects in the trenches with such crystal defects.\nRecently techniques in the fields of quantum information and quantum computation attract a great deal of attention as described above, and possibilities of applications of quantum dots attract attention. However, for further fundamental studies and application development of quantum dots, various technical barriers which have to be overcome are present.\nFor example, quantum dots have small sizes, and the quantum dots self-assembled by using S-K mode are distributed at random. No method that makes the self-assembled quantum dots distributed at random electrically accessible has been so far proposed.\nOn the other hand, conductor pads of an about 100 nm-size can be formed by using processes, such as electron beam lithography, reactive ion etching, etc., which can make submicron processing. Accordingly, in a case that the density of quantum dots is very low, it will be possible to form an electrode above a single quantum dot by using such process. That is, electrodes are formed suitably above the regions where quantum dots are expected to have been formed, by using a process which can make submicron processing. Then, whether or not the quantum dots are present below the electrodes is checked. The electrode formation and the following check are repeated many times, whereby the presence of a single quantum dot below the electrode could be often found. However, it is almost impossible to capture a single quantum dot in a case that the density of the quantum dots is high. Such method depending on the probability is inefficient to fabricate devices.\nJapanese published unexamined patent application No. Hei 07-297381 (1995) proposes a method of electrically accessing quantum dots by means of electrodes in the form of a fine probe shape. However, in this case, in order to arrange electrodes above quantum dots, a special process for forming the quantum dots in an array or others is necessary. In Patent Reference 1, electric fields of the probe-shaped electrodes will be distributed far more widely than a size of the quantum dots. Accordingly, the electric fields of the electrode arranged above a certain quantum dot will influence quantum dots adjacent to said certain quantum dot.\nAs described above, the conventional techniques have established no method of making accurate electric access to each of quantum dots which have self-assembled at random.\nIt is very significant in various aspects of the fundamental studies, application developments, etc. of quantum dots to realize the accurate electric access to each of self-assembled quantum dots.\nAn object of the present invention is to provide a method for forming a quantum dot which can form a quantum dot of good-quality in a prescribed position and in a prescribed size.\nAnother object of the present invention is to provide a quantum semiconductor device which can make accurate electric access to a quantum dot and a method for fabricating the quantum semiconductor device."}
{"text": "Nowadays, a conventional antenna cannot be used in an electronic device with a full metal shell, since a signal of the antenna may be shielded and thus the antenna cannot work normally. Furthermore, the conventional antenna of the electronic device cannot cover WWAN (Wireless Wide Area Network)."}
{"text": "Reclosable plastic bags are used for many purposes, including transporting and storing food products and other materials. Reclosable bags are often preferred over other available bags because of numerous advantages, including that they are easy to close and reopen for access without compromise of integrity when closed.\nTypically, such bags are reclosable using a zip-action, locking reclosable closure. Unfortunately, when using such a bag for a first time, there is currently no way of knowing if the bag is sanitary, has been used, tampered with or otherwise contaminated. This is particularly important if the contents to be stored in the bag are perishable or easily contaminated. Exemplary of such is breast milk.\nLikewise, currently available bags are typically difficult to fill (E.G., with a liquid substance) without obstruction from the lips and/or spilling along the sides of the bag. At least one reason it is difficult to fill currently available bags is that the edges are typically heat-sealed along their entire length from the bottom of the bag to the top of the bag. These heat-sealed edges typically make it somewhat difficult to pull the lips apart, as well as impede the ability of the bag to remain open during filling.\nYet another disadvantage of using currently available bags to store liquids is that the zipper may not remain sealed. This is particularly true when, for instance, the bag lies on its side—the bag thus situated is more likely to seep some of the liquid contents.\nThe present invention avoids these and other difficulties by providing a tamper evident reclosable bag that ensures that the bag is being used for the first time. The present invention further provides a bag that is easier to open, easier to fill, seals better, and has less leakage than currently available bags."}
{"text": "The present invention relates to a solid state image sensor comprising a plurality of pixels formed by static induction transistors having both a photoelectric converting function and a switching function.\nHeretofore, there have been proposed various solid state image sensors comprising a charge transfer device such as CCD and BBD, and MOS transistors. However, such known image sensors have drawbacks in that electric charges are leaked during the charge transfer and the photodetection sensitivity is low. In order to avoid such drawbacks, there has been proposed a solid stage image sensor comprising static induction transistors (hereinafter abbreviated as SIT). For instance, in an European Patent Application No. 83900059.3 (Publication No. 96725), there is described a solid state image sensor comprising SITs each of which serve as a photodetecting element and a switching element. The inventor of the instant application has also proposed in U.S. patent application Ser. No. 647,169 filed on Sept. 4, 1984, a solid state image sensor comprising SITs., which solid state image sensor has a higher sensitivity and is easily manufactured. This solid state image sensor can read out an image signal while stored charges remain undestroyed and can have a high sensitivity by prolonging the time period of accumulating photocarriers.\nFIGS. 1A and 1B show an embodiment of the solid state image sensor comprising SITs as proposed in the above mentioned Patent Application. A SIT shown in FIG. 1A comprises an n.sup.+ substrate 1 which serves as a drain of the SIT, a lightly doped n.sup.- epitaxial layer 2 which is grown on the substrate 1 and serves as a channel, an n.sup.+ source region 3 and p.sup.+ gate regions 4 formed in the epitaxial layer 2 by means of, for example, a thermal diffusion, source electrode connected to the source region 3 and gate electrodes 6 arranged above the gate regions 4 via an insulating film 5 such as SiO.sub.2 to form a gate capacitor 7. The SIT is isolated from adjacent SITs by means of an isolation region 9 formed by a burried insulating substance. A number of SITs are arranged in a matrix on the same substrate.\nFIG. 1B is a circuit diagram illustrating a whole construction of the solid state image sensor comprising normally-on type SITs shown in FIG. 1A. Drains (substrate) of SITs 10-11 to 10-mn forming pixels arranged in a matrix are commonly connected to ground, and to gate electrodes of SITs 10-11 to 10-1n; . . .; 10-m1 to 10-mn arranged in the X direction, i.e. in a row, are connected row lines 11-1 to 11-m, respectively, which lines are then connected to a vertical scanning circuit 16 to receive row selection signals .phi.G1 to .phi.Gm. Source electrodes of SITs 10-11 to 10-m1; . . .; 10-1n to 10-mn arranged in the Y direction, i.e. in a column are connected to column lines 12-1 to 12-n, respectively, whose first ends are connected to a video voltage source Vs via respective column selection transistors 13-1 . . . 13-n, a common video line 14 and a load resistor 15. To gates of the column selection transistors 14-1 to 14-n are applied column selection signals .phi.S1 to .phi.Sn, respectively from a horizontal scanning circuit 17.\nIn the SIT having the construction explained above, when a light input is given, electron-hole pairs are induced in the channel region 2 and a gate depletion region. Electrons are conducted away into the drain 1 connected to ground, and holes are stored in the signal storing gate region 4 and thus the gate capacitance 7 connected thereto is charged to vary a gate potential by .alpha.V.sub.G. Now it is assumed that an amount of electrostatic charges stored in the gate region 4 is represented by Q.sub.L and the capacitance of the gate capacitor 7 is denoted by C.sub.G, then .alpha.V.sub.G =Q.sub.L /C.sub.G is obtained. After a certain accumulation time, when a readout pulse V.sub..phi.G is applied to the gate terminal 8, the gate potential is changed into a sum of V.sub..phi.G and .alpha.V.sub.G. Then, a potential difference between the signal storing gate region 4 and source region 3 is lowered to decrease the depletion layer and a drain current having an amplitude corresponding to the light input begins to flow between the source and drain. Due to the amplifying function of SIT, the drain current is multiplied by the amplification factor of the SIT. It should be noted that even if the source and drain of the SIT are exhanged, a similar operation to that explained above may be effected.\nFIGS. 2A to 2F illustrate waveforms of the signals applied to the row lines 11-1 to 11-m and column selection transistors 14-1 to 14-n. As illustrated in FIGS. 2A to 2F, in this solid state image sensor, the successive pixels are readout by successively selecting the row lines 11-1 to 11-m as well as by successively selecting the column lines 13-1 to 13-n while a row line has been selected. Each of the signals applied to the row lines has a lower voltage V.sub..phi.G, and a higher voltage V.sub..phi.R. During the line scanning period t.sub.H the row selection signal has the voltage V.sub..phi.G and during a horizontal blanking period t.sub.BL the row selection signal assumes the voltage V.sub..phi.R. Each of the column selection signals .phi.S1, .phi.S2 . . . applied to the gates of the column selection transistors has a low level for cutting off the transistor and a high level for making the transistor conductive.\nWhen the vertical scanning signal .phi.G1 is increased to the high voltage V.sub..phi.G, SITs connected to the first row line 11-1 are selected and the SITs are successively read out while the column selection transistors 13-1, 13-2 . . . 13-n are successively made conductive by means of the horizontal scanning signals .phi.S1, .phi.S2 . . . .phi.Sn. In this manner, a video signal of one line is derived on the video line 14. Then, the SITs in the first row are simultaneously reset when the vertical scanning signal .phi.G1 is increased to the higher level V.sub..phi.R.\nNext, when the vertical scanning signal .phi.G.sub.2 is changed to the voltage V.sub..phi.G, all SITs 10-21, 10-22 . . . .phi.-2n connected to the second row line 11-2 are selected and are successively readout by means of the horizontal scanning signals .phi.S1, .phi.S2 . . . .phi.Sn.\nIn this manner, successive pixels are readout to derive a video signal of one frame.\nThe inventor has found that the known solid state image sensor mentioned above has the following drawbacks. In the image sensor, the accumulation of photocarriers induced by the incident light is effected in the channel region 2 and gate depletion layer in each SIT. The accumulating regions correspond to portions in the channel region which situated between the gate regions and isolating region, and underneath the gate regions. It is apparent that the portions between the gate regions and isolating region are liable to be small as the dimension of pixels is made smaller. Then, a major portion which serves to accumulate the photocarriers would be the portion in the channel region situated underneath the gate region 4, and thus, the incident light has to transmit through the gate regions 4. In order to obtain a large amplification of the SIT, a depth Xj of the gate region 4 has to be relatively thick such as 2 to 4 .mu.m. Therefore, the incident light, particularly components of short wavelengths (blue light) is abosrbed in the gate region 4 and cannot penetrate to the portion of channel region underneath the gate region.\nFIG. 3 is a graph showing a light absorption coefficient .alpha. of silicon in dependence upon wavelength .lambda. of incident light. If it is assumed that .lambda.=0.4 .mu.m and the thickness of the P.sup.+ gate region is 3 .mu.m, then .alpha.=6.times.10.sup.-4 cm.sup.-1, and a light transmittivity becomes a very small value such as e.sup.-6.times.10.spsp.4.sup..times.3.times.10.spsp.4 =1.5.times.10.sup.-8. In this manner, the known solid state image sensor has a very low sensitivity especially for blue light and thus a color camera comprising the solid state image sensor having a high sensitivity over a whole sepectrum of visible light can not be realized."}
{"text": "The present invention relates to a teletext converter which is external to a television receiver (\"set-top\" converter), and more particularly, to such converters having a \"transparent mode\".\nTeletext is a service that broadcasts digital information within the vertical blanking interval (VBI) of a standard broadcast television signal and at the receiver presents the information as digitally generated text and/or pictures on a display screen. In the most commonly considered mode of operation, the standard TV picture is completely replaced by the digitally generated picture. \"Transparent mode\" refers to the situation in which some region of the screen, rather than being defined by the teletext signal, is occupied by the regular video signal. One example is a regular video picture having a teletext generated sub-title or caption. Ideally, the teletext converter, which converts the digital signals into video signals, is located within the receiver. However, almost all existing receivers do not have such an internal converter. Thus an external set-top converter is used if teletext signal reception is desired on most existing sets. A set-top converter receives the antenna signal, converts the digital teletext signal into a video signal, and then provides an output signal for connection to the receivers antenna input terminal.\nOne prior art method of producing a transparent mode set-top converter is to decode the regular video to baseband red, green and blue signals, combine these signals with R, G, and B signals derived from the decoded teletext, re-encode to composite NTSC, and then modulate with the composite NTSC signal an output carrier having a selected channel frequency. This approach suffers the deficiencies of high cost and significant signal quality degradation.\nIn particular, the decoding of composite video to baseband color signals, involving the separation of luminance and chrominance, band-limiting these signals and various other distortion producing processes degrades the fidelity of the standard TV picture produced after re-encoding and finally decoding again in the TV receiver.\nIt is therefore desirable to provide a set-top teletext decoder with transparent mode that does not cause signal degradation and has accurate color fidelity."}
{"text": "Recently, the technique for reducing mounting costs is studied and developed by use of a coupling system called a direct optical coupling (butt joint) system. In the coupling system, an optical fiber and optical semiconductor element are arranged in close positions to face each other and optically coupled without using a lens. When the direct optical coupling system is used, it is important to arrange an optical semiconductor element and optical fiber in close positions and prevent light from reaching an unwanted portion. Therefore, there is proposed a method for directly forming an electrical interconnection on the main surface of an optical transmission line holding member (a so-called optical fiber ferrule) that holds the optical fiber and the like and mounting an optical semiconductor element on the main surface to set the light receiving/transmitting region thereof in opposition to the optical fiber (Jpn. Pat. Appln. KOKAI Publication No. 2001-159724).\nIn this type of method, it is desired to fill resin in between the optical semiconductor element and the optical transmission line so as to suppress reflection light caused by discontinuity of the refractive index. Specifically, resin having a role as a refractive index matching material and an under-filling function for protection of the electrical connecting portion of the optical element is filled in between the optical semiconductor element and the end face of the optical transmission line. Further, resin used to mechanically fix the optical transmission line on a ferrule is filled.\nAs a method for forming the above resins, a method for performing an under-filling process before fixing the optical transmission line and a method for fixing the optical transmission line and then performing an under-filling process are considered. When the optical transmission line is first fixed, there occurs a possibility that fixing resin flows into holding holes for the optical transmission line to reach the end face on the optical element side. Further, since the fixing resin fills in the holding holes, the holes via which air escapes when under-fill resin is injected are closed or blocked. This means that the holes via which a solvent produced at the curing time of the under-fill resin is extracted are closed. Therefore, there occurs a possibility that bubbles remain in the resin at the curing time of the under-fill resin and this may prevent the optical coupling. When the under-filling process is first performed, there occurs a problem that high mechanical precision cannot be attained."}
{"text": "Coplanar waveguides (CPW) generally comprise a conductive line placed between two ground lines, and allow components such as antennas operating at frequencies of, for example, about several tens of gigahertz, to be electrically connected. When a high-frequency signal is passed through such a waveguide, losses can lead to the signal being attenuated. These losses are in particular due to the resistance and inductance of the line, to capacitive coupling to the substrate, and to conductance allowing ground return.\nIn order to reduce these losses, it has been suggested, by Gamble, “Low-loss CPW lines on surface stabilized high-resistivity silicon,” IEEE Microwave and Guided Wave Letters, Vol. 9, No. 10, October 1999, that these coplanar waveguides be fabricated using metal lines on high-resistivity silicon substrates. Thus low conductance is obtained between the conductive line and the ground lines. The diffusion of metals, for example gold or copper, into the silicon, and the possibility that charge trapped in a passivation layer placed in contact with the conductive lines will induce a conductive channel in the substrate, between the signal and ground, are the drawbacks of this solution.\nSuperposing the conductive lines of the coplanar waveguides on insulating lines, for example silicon-dioxide (SiO2) lines, on a high-resistivity silicon substrate has also been suggested. This solution has the drawback of forming capacitive coupling between the conductive lines and the substrate. Moreover, charge trapped in a passivation layer, placed in contact with the conductive lines, can also induce a conductive channel in the high-resistivity silicon substrate, between the regions under the conductive lines.\nIt has also been suggested that the conductive lines of the coplanar waveguides be formed on an insulating layer, placed on a polycrystalline-silicon layer, itself placed on a high-resistivity silicon substrate. This solution allows current flowing under the insulating layer to be captured in the polycrystalline-silicon layer.\nThe reader may also refer to Ojefors, “Micromachined loop antennas on low resistivity silicon,” IEEE Transactions on Antennas and Propagation, 2006, Vol. 54, No. 12, pages 3593-3601, describing the formation of trenches in the silicon substrate under the oxide layer and located substantially between the conductive lines of an antenna and a ground plane. This solution in particular has the drawback of being complicated to implement.\nFinally, Wojnowski, “Package Trends for Today's and Future mm-Wave Applications,” Infineon, 38th European Microwave Conference, 2008, describes the formation of a cavity matrix under the coplanar waveguides. This solution does not sufficiently limit losses, in particular because of parasitic conductivity in the silicon. The formation of the cavities described in this document also has the drawback of being implemented before metal lines have been formed on the frontside of the integrated circuits.\nMoreover, it may be necessary to use through-silicon vias in integrated circuits, for example in order to form three-dimensional stacks or in order to form electrical contacts on a side opposed to that comprising the metal lines of the integrated circuit."}
{"text": "1. Field of the Invention\nThe present invention relates to controllers for servomechanisms, and in particular, to a digital servomechanism controller that includes both a \"dumb\" servomechanism controller and a neural network based servomechanism controller working in conjunction with one another.\n2. Description of the Related Art\nComputer disk drive servomechanisms are electromechanical devices which require control systems for controlling their rotational velocity and the positioning of their read, write and position sensing heads. Such control systems, commonly referred to as servo controllers, are well known in the art. As illustrated in FIG. 1, a servo controller and the servomechanism to be controlled are typically interconnected in a classical feedback control system design.\nThe servo controller provides a servo control signal to the \"plant,\" i.e. the disk drive servomechanism, for controlling the performance thereof, e.g. to rotate the storage media and position the read, write and position sensing heads. The plant generates feedback signals which represent actual performance parameters of the plant, e.g. the actual velocity of the storage media and the actual position of the heads. Each actual performance signal is combined with a reference signal which represents the desired performance of the plant. This produces an \"error\" signal (representing the difference between the desired and actual performance signals) which is used to drive the servo controller.\nThe servo controller adjusts its servo control signal output in accordance with its error signal input. As the performance of the plant changes in accordance with its servo control signal input, the actual performance feedback signal more closely resembles the reference signal, thereby causing the error signal to decrease. Once the error signal has decreased below a predetermined threshold value, the servo controller ceases to modify its servo control signal output, and instead simply maintains its servo control signal output until such time as the error signal once again rises above the predetermined threshold.\nThe simplest servo controller is a \"dumb\" (e.g. non-programmable, or with a fixed or static program) analog network which accepts an analog error signal and provides an analog servo control signal to the plant in accordance with a substantially fixed design or algorithm. In turn, the actual performance feedback signal and reference signal are analog signals as well. Typical control characteristics provided by analog servo controllers are referred to as \"proportional,\" \"integral,\" \"derivative\" and \"proportional-integral-derivative\" (\"PID\").\nA \"proportional\" type of servo controller produces a servo control signal which is proportional to the error signal and can be described as follows:\nU=K.sub.p E PA1 U=servo control signal PA1 E=error signal PA1 K.sub.p =proportional feedback coefficient PA1 dt=time differential\nwhere:\nAn \"integral\" type of servo controller produces a servo control signal based upon the integral of the error signal and can be described as follows: ##EQU1## where: T.sub.I =integral feedback coefficient\nA \"derivative\" type of servo controller produces a servo control signal which is based upon the derivative of the error signal and can be described as follows: ##EQU2##\nA \"PID\" type of servo controller produces a servo control signal which can be proportional to the error signal, while simultaneously containing components based upon the integral or derivative of the error signal, or both. Such a servo control signal can be described as follows: ##EQU3##\nThe relative simplicity of analog servo controllers, however, renders them inadequate for controlling today's high performance disk drive systems having increasingly denser storage capabilities. This has resulted in the widespread use of digital servo controllers.\nDigital servo controllers, particularly those using microprocessors, allow for more sophisticated control characteristics. By properly programming its microprocessor, a digital servo controller can also be a \"proportional,\" \"integral,\" \"derivative\" or \"PID\" type of servo controller. But, the flexibility afforded by the programmability of the digital servo controller allows more precise control to be exercised.\nHowever, digital servo controllers have limitations of their own. Even with their programmability, their design, whether with respect to hardware or software, is essentially fixed. In other words, the control characteristics are fixed in the sense that they cannot anticipate and compensate for the varying control environment and system nonlinearities which inevitably exist but cannot be precisely predicted or modeled.\nAs will be appreciated, variances in the control environment can arise and system nonlinearities can be introduced in a number of ways. For example, variances in the control environment can arise due to changes in the operating characteristics or tolerances of the hardware components used to construct the system. Such changes can be induced by factors such as changes in temperature, humidity, shock, vibration, power supply voltages, etc. System nonlinearities can be introduced by system hardware components which vary randomly from their nominal rated values, and yet are still within their rated tolerances.\nTo overcome the limitations of a simple digital servo controller, \"adaptive controllers\" have been developed. As illustrated in FIG. 2, an adaptive controller is used in conjunction with a digital servo controller. Two error signals are produced, error signals \"A\" and \"B,\" for use by the digital servo controller and adaptive controller, respectively.\nThe actual performance signal produced by the plant, in addition to being compared with the reference signal directly, is also compared to the output signal of the \"reference model.\" The reference model is a preprogrammed digital circuit, or software model, which receives the reference signal and produces an output signal based upon a desired copy of the servomechanism system. In other words, the reference model is a mathematical model based upon historical data regarding the control environment and system nonlinearities of the system, including the plant and servo controller. This model is then used to generate a desired model output signal.\nThe result of this comparison between the actual performance signal and the modified reference signal, i.e. error signal B, is used to drive the adaptive controller. In accordance therewith, the adaptive controller provides a signal to the servo controller which is used to adjust the parameters of the control characteristics used by the digital servo controller in producing the servo control signal for the plant.\nFor example, the signal from the adaptive controller can be used to adjust the \"proportional,\" \"integral\" or \"derivative\" constants or coefficients used by the servo controller in generating the servo control signal. By using the adaptive controller in this manner, the otherwise fixed design or algorithm of the servo controller can be adapted in accordance with the reference model to more closely compensate for the varying control environment and system nonlinearities. In other words, using the adaptive controller in this way allows the overall system operation to track the operation model stored as the reference model.\nHowever, adaptive controllers tend to be quite complex and typically require large amounts of computation time. Furthermore, mathematical models (e.g. for the reference model) which give good representations of the servomechanism control environment and the disk drive system are very difficult, if not impossible, to achieve.\nThe overall operating environment sought to be controlled is inherently nonlinear. Therefore, the requisite mathematical models should ideally also be nonlinear so as to more closely approximate that environment. However, such nonlinear models provide ad hoc control solutions and are unique for each application to avoid introducing excessive errors of their own. Further, they are generally quite complex and difficult to redesign for use in other applications.\nLinearized models, i.e. models which are based upon linear mathematical models, have been used with some success. However, while linear models are simpler and may be more easily redesigned for use in other applications, their performance outside of the relatively narrow regions for which they are optimized has been disappointing.\nTo meet this need for servo controllers which can adaptively control inherently nonlinear servomechanisms in a varying control environment, intelligent servo controllers have been proposed. An advantage of intelligent servo controllers is their ability to provide more reliable generalized solutions for controlling servomechanisms over wider ranges of operating conditions and uncertainties. This is in contrast to the ad hoc control solutions provided by adaptive controllers using nonlinear models which tend to be application specific.\nOne type of proposed intelligent servo controller uses artificial intelligence. Another proposed type, which offers better performance, is referred to as \"fuzzy.\" However, the artificial intelligence or \"fuzzy\" controllers are still computation intensive and based upon rules or algorithms. Therefore, their generalization capabilities are poor. Thus, they have not been widely accepted.\nA third type of proposed intelligent servo controller, which has offered the best performance so far, is that which uses a neural network. Its advantage is that a neural network is not reliant upon algorithms and is capable of adapting to virtually all parameter or system variations. Therefore, it has better generalization capabilities.\nUsing a neural network servo controller to control a disk drive servomechanism requires that the neural network servo controller be coupled to the plant in two distinct configurations, as shown in FIGS. 3A-3B. As shown in FIG. 3A, the initial configuration establishes the general learning phase for the neural network servo controller. The neural network servo controller receives the actual performance signal from the plant and adjusts its adaptive weights in accordance with the error signal produced by the comparison of the reference signal with the output signal of the neural network servo controller. In this configuration, the neural network servo controller learns the inverse model of the plant broadly.\nThe second configuration, shown in FIG. 3B, is the specialized learning phase for the neural network servo controller. The reference signal drives the neural network servo controller directly, and is also compared to the actual performance signal output of the plant to produce the error signal. The error signal is then used to further train the neural network. Therefore, in this configuration, the neural network servo controller controls the plant via the servo control signal, while simultaneously learning and adapting to variations in the control environment and system nonlinearities.\nHowever, the use of a neural network based servo controller has limitations with respect to its practicality. The general learning phase, as illustrated in FIG. 3A, is very broad. A large number of hardware neurons, or software based calculations, are needed to take advantage of the higher resolution and accuracy capabilities provided by the neural network. Furthermore, to achieve maximal generalized learning, rigorous learning in the form of using many sets of input patterns, i.e. reference signals, is required.\nTherefore, it would be desirable to have a disk drive servo controller which provides the resolution, accuracy and adaptability of a neural network, but which does not require an impractically large number of hardware neurons or software based calculations, nor require extensive input learning patterns."}
{"text": "The method and the apparatus of the present invention are applicable to a hot-dip coating of a ferrous base metal strip with zinc, zinc alloys, aluminum, aluminum alloys, terne, lead and those coating metals or coating metal alloys which have oxide forming characteristics such that optimum finishing cannot be accomplished by conventional jet practice or by conventional exit rolls. While not intended to be so limited, for purposes of an exemplary showing the method of the present invention will be described as applied to galvanizing. The method can be practiced on various types of galvanizing lines. For example, the method of the present invention is applicable to fluxless, hot-dip metallic coating of ferrous base metal strip where it is necessary to subject the strip surfaces to a preliminary treatment which provides the strip with oxide free surfaces and preferably brings the strip to a temperature approximating that of the molten zinc or zinc alloy coating bath at the time the strip is caused to pass beneath the surface thereof. One of the principal types of anneal-in-line, fluxless, preliminary treatments is the so-called Sendzimir process or oxidation-reduction practice disclosed in U.S. Pat. Nos. 2,110,893 and 2,197,622. Another anneal-in-line, fluxless, preliminary treatment in common use is the so-called Selas process or high intensity direct fired furnace practice disclosed in U.S. Pat. No. 3,320,085.\nIn the Selas process the ferrous base metal strip is passed through a direct fired preheat furnace section. The strip is heated by direct combustion of fuel and air producing gaseous products of combustion containing at least about 3% combustibles in the form of carbon monoxide and hydrogen. The strip reaches a temperature of about 535.degree. C. to about 760.degree. C. while maintaining bright surfaces completely free of oxidation. The strip is then passed into a reducing section which is in sealed relation to the preheat section and which contains a hydrogen and nitrogen atmosphere, wherein it may be further heated by radiant tubes to from about 650.degree. C. to about 925.degree. C. and thereafter cooled approximately to the molten coated metal bath temperature. The strip is then led beneath the bath surface while surrounded by the protective atmosphere.\nOther related pretreatment techniques are taught in U.S. Pat. No. Re 29, 726--U.S. Pat. Nos. 3,837,790--4,123,291--4,123,292 and 4,140,552. The above mentioned prior art patents constitute non-limiting examples of fluxless, continuous galvanizing processes to which the method of the present invention is applicable. When such conventional strip preparation techniques as are taught in the above mentioned prior art are used, it is necessary that the base metal strip be maintained in a protective atmosphere at least until it passes beneath the surface of the bath of molten zinc or zinc alloy.\nSuch a protective atmosphere is not a requirement when flux or chemical strip preparation techniques of the type taught in U.S. Pat. Nos. 2,824,020 and 2,824,021 are employed. Briefly, when such chemical strip preparation techniques are used, the ferrous base metal strip is caused to pass through a flux bath and through means to assure the proper thickness of the flux coating on the strip. The ferrous base metal strip is then conducted through a heating chamber wherein the strip is heated to evaporate the water in the flux solution. Thereafter, the ferrous base metal strip is further heated to raise it to that temperature approaching the maximum temperature of stability of the flux coating on the strip. The strip is then caused to pass beneath the surface of the bath of molten zinc or zinc alloy so as to be coated. The method of the present invention is equally applicable to galvanizing lines utilizing such flux or chemical pretreatment systems.\nFrom the above it will be evident that the method of the present invention is not limited to the use of any particular pretreatment of the ferrous base metal strip in the galvanizing line and the terms \"pretreatment\" or \"pretreated\" (as used herein and in the claims with reference to the ferrous base metal strip) are to be interpreted broadly to include any of the conventional pretreatment systems exemplified by the above noted prior art. In general, these terms refer to any appropriate pretreatment technique, the result of which is such that, during the actual coating step wherein the ferrous base metal strip passes through the molten bath of zinc or zinc alloy, it will be at or will achieve the proper coating temperature and its surface will be oxide-free.\nIn conventional continuous hot dip galvanizing, it is usual to cause the two-side coated strip to exit the molten coating metal bath into the ambient atmosphere. The most widely used finishing and coating weight control technique is to direct the coated strip between jet knives or nozzles which cause a blast of air or steam to impinge upon both sides of the coated strip, returning excess coating metal to the bath. This finishing technique, however, has a number of definite drawbacks. Some of these drawbacks include coating ripples, oxide curtains and, bath surface oxide problems, including the necessity for top skimming. All of these defects are real, but are reasonably controllable by known methods.\nU.S. Pat. Nos. 4,107,357 and 4,114,563 and German Pat. No. 2,656,524 are exemplary of patents teaching methods for coating one side only of a ferrous base metal strip. In the practice of these processes, the coated strip after contact with the coating bath is maintained in a protective, non-oxidizing atmosphere and is jet finished with nitrogen or non-oxidizing gas. However, the primary purpose of these steps is to prevent oxidation of that side of the ferrous base metal strip not coated or, if the uncoated side has an oxide film thereon, to prevent adherence of the coating metal to the oxide film. Nevertheless, U.S. Pat. No. 4,114,563 mentions in passing that the above noted defects are controlled or eliminated when finishing in a protective atmosphere.\nA major problem area encountered with conventional jet finishing is that of coating control at the edges of the strip. One edge problem is that of zinc coating thickness over a narrow band immediately adjacent each edge of the coated strip. The coating thickness of these bands is greater than the coating thickness over the rest of the strip width. If this coating thickness differential is great enough, edge buildup or spooling will occur when the continuous strip is coiled under tension.\nEdge berries (small balls of oxide) which are attached to the strip edge and are pulled through the jet blast constitute another troublesome problem. Furthermore, an edge defect commonly known as \"feathered oxide\" occurs during low speed jet finishing. Feathered oxide is characterized by discontinuous patches of heavy coating metal oxide which pull through the jet blast. They appear much like feathers which extend inwardly from the strip edges with the tips thereof pointing toward the center of the strip.\nMany methods have been used by prior art workers to reduce build-up and oxide control problems at the strip edges. Tapered jet nozzle slot openings are commonly used where the slot opening of the jet finishing nozzle continuously increases in width from the center of the jet nozzle to its ends. Such a contoured jet finishing nozzle is taught in U.S. Pat. No. 4,137,347.\nOther methods to control edge coating include curving the jet nozzles so that the nozzle is closer to the strip at the strip edges than at the strip center. Also, vanes or nozzle extensions have been used at the strip edges to bring the nozzle closer to the edges than to the center of the strip. Still other methods include the use of shutters and auxiliary jets, both internal and external to the main jet nozzles, to alter the jet wiping force at the strip edges as compared to the jet wiping force at the strip center.\nAll prior art methods fall short of producing optimum edge control with a minumum of operator attention, maximum coating metal economy, and proper edge control over a wide range of strip widths and line speeds.\nEdge build up in a one-side coating process differs from that encountered in a conventional two-side coating process. This is true because, when jet finishing is employed, it is normally performed with a single jet nozzle, preferably with a back-up roll on the uncoated side of the strip. As a result, edge build-up is less severe, since the back-up roll serves as an extension of the strip edges.\nYet another major hot-dip zinc coating defect is commonly referred to as \"spangle relief\". Spangle relief has two aspects. One is the variation of surface profile (zinc thickness) across the zinc crystal from one boundary to the opposite boundary. The other is a depressed spangle boundary which surrounds each spangle or crystal. Spangle relief can be reduced by such methods as purposely causing part of the zinc coating to alloy with the ferrous base metal, or by antimony additions to the zinc bath. However, none of these methods is entirely satisfactory.\nAs a result, many methods have been developed to suppress spangle formation. That is, to minimize final spangle size to such an extent that the spangles are hardly visible to the naked eye. For example, U.S. Pat. Nos. 3,322,558; 3,379,557 and 3,756,844 teach spangle minimizing methods. Most methods involve spraying water or water solutions against the molten coating to quench the coating and create many nucleation sites. However, the results achieved by these spangle minimizing methods are not always consistent.\nWhile spangle relief can be reduced or masked by temper (skin-pass) rolling such temper rolling causes these defects to become imprinted in the base metal. As a result of this non-uniform cold working of the base metal, the defects may reappear when critical surface items such as automotive body parts and appliance parts are stamped or formed. These defects may not be completely masked after the parts are painted.\nCoating irregularity problems constituted a major reason why prior art workers have recently turned their attention to a one-side coated product (as exemplified by U.S. Pat. No. 4,082,868) wherein the uncoated side is the side to be painted on critical surface items made therefrom, even at the sacrifice of corrosion resistance on the non-galvanized side.\nYet another problem area encountered with conventional jet finishing involves coating weights and line speeds. The viscous interaction between the coating metal and the strip is proportional to strip speed. At slow speeds, the prior art was faced with the problem of ripple formation. To combat this, it was found that minimizing jet nozzle to strip distance and reducing the jet finishing pressure will minimize the prominence of coating ripples. However, low jet finishing pressure and close positioning of the jet finishing nozzles at the same time create an edge build-up problem. Prior art workers therefore have had to adjust the parameters to control edge build-up and ripples and this has necessitated higher line speeds. As an example, it has been common practice to use conventional jet finishing only with strip speeds above about 100 feet per minute (30 meters per minute) to produce commercial class coating weight (ASTM A525, G90). Edge build-up problems on G-90 coating commonly occur at speeds below about 150 feet per minute (45 meters per minute). Heavier coatings are even more difficult to control edge-to-edge uniformly.\nAnother essential practice in most prior art jet finishing operations is to position the jet nozzles virtually directly opposed, such that the jet streams are in direct interference beyond the edges of the strip. This interference results in extremely high and objectionable noise levels. Vertically offset jet nozzles cause a wrap-around effect with a bead of heavy coating metal along the edge on the side opposite the last nozzle operating on the strip. In addition to the noise problem and the need for precise adjustment of the jet nozzles by the operator, opposed operation can result in coating metal splatter being blown off of the strip edge by one nozzle and into the nozzle opening of the opposed nozzle.\nPrior art workers have hitherto jet finished hot-dipped, two-side coated galvanized and aluminized strip with nitrogen. Such jet finishing, however, has been performed in an ambient atmosphere. In jet finishing, less nitrogen is required when nitrogen is used than air when air is used. However, the results achieved by such finishing are more nearly like those achieved in jet finishing with air in an ambient atmosphere than like the results achieved by the present method.\nU.S. Pat. Nos. 3,505,042 and 3,505,043 teach hot dip coating of a ferrous metal with magnesium-zinc alloy and magnesium-aluminum-zinc alloy respectively. The first of these references utilized coating rolls for coating control and nitrogen gas solely for cooling after finishing between exit rolls. The second of these references does not disclose specific finishing means, but does enclose the coating apparatus in an atmosphere of 10% hydrogen with the balance nitrogen. The reference also mentions blowing nitrogen gas over the coated strip being withdrawn from the bath to rapidly cool the coating and to control coating weight, the nitrogen gas having a temperature below 100.degree. F. (40.degree. C.) and preferably about 50.degree. F. (10.degree. C.).\nThe present invention is based upon the discovery that if, in a conventional, continuous, hot-dip, two-side galvanizing process, the coating metal as it exits the coating bath is surrounded by an enclosure in which a substantially oxygen-free atmosphere is maintained and if within the enclosure the coated strip is jet finished with non-oxidizing or inert gas, the finishing problems encountered with conventional finishing methods are markedly reduced or eliminated. The finishing method of the present invention enables the practice of any of the above noted spangle minimizing techniques, greatly enhancing the results achieved and the uniformity of spangle minimization throughout the width of the strip. Furthermore, the markedly improved results are consistent.\nThe most significant aspect of the present invention, however, is the discovery that all coating control problems at the strip edges are completely eliminated with the substantial exclusion of oxygen from the finishing process. Minimum operating speeds are no longer limited by edge build-up problems, but rather only by the desired coating weight relative to the amount of coating metal naturally pulled up by the strip to the finishing jet nozzles. It has been found, for example, that excellent quality G-90 coatings can be produced without difficulty at speeds as low as 30 feet per minute (9 meters per minute).\nCoating control at the strip edges is a very real problem in a two-side coating process where jet finishing cannot be accomplished against a back-up roll as in one-side coating processes. In the practice of the teachings of the present invention, heavy coating at the strip edges does not occur. Jet nozzles can be vertically offset, eliminating the need for precise positioning, greatly reducing noise, and eliminating the hazard of zinc splatter plugging opposed nozzles. Jet nozzle design may be simplified to utilize a slot-like nozzle opening of uniform width throughout its length, rendering unnecessary the multitude of special jet nozzle designs, methods and accessories which have been used for controlling edge build-up. Superior uniform coating thickness results edge-to-edge for all coating weights because the center profile need no longer be distorted to compensate for heavy edges. As a result, coils can be wound under high tension without spooling. Furthermore, with uniform coating weight edge-to-edge, less coating metal is used. The greatly reduced noise, by virtue of vertically offsetting the jet nozzles in an enclosure, represents a major step in safety and environment improvement. The simplified jet nozzle design results in more uniform coating, more uniform cooling and a flatter strip.\nHeretofore, in the practice of conventional two-side jet finishing, neither the mechanism for ripple formation nor that causing edge build-up problems was completely understood. In the process of conventional jet finishing, a pneumatic dam effect is created whereby the desired amount of coating metal is metered through the jet barrier to form the finished coating. At this metering point the excess coating metal pulled up with the strip, beyond that required for the finished coating, is returned to the coating bath. This process is described in detail in U.S. Pat. No. 4,078,103.\nWhile applicants do not wish to be bound by theory, it appears as a result of the present invention that the coating ripples and heavy edge coating in conventional jet finishing are caused entirely by coating metal oxide. At some point in the jet interaction region, probably just above the point of zero surface velocity , fresh (unoxidized) coating metal is being exposed and, as it is exposed, it immediately forms a very light oxide skin. The continuity of flow or distribution of this very light oxide skin onto the finished coating determines the occurrence of coating ripples. In conventional practice, the jet periodically restrains the oxide film. The film builds up until the jet no longer can restrain it. At this time a segment of relatively heavy oxide breaks off and passes with the finished coating. The segment as it passes carries with it coating beneath which is heavier than that which is metered on when the oxide film is restrained. This process is repeated many times each second as ripples are formed.\nA similar mechanism is believed to be operable in creating heavy coating metal along the strip edges. However, at the edges, geometry becomes an additional important factor in that there is no wiping force directed against the edge surfaces. Relatively heavy oxide is permitted to pass through the jet interaction area more or less continuously carrying with it heavy coating beneath. This oxide envelope around each strip edge surface is the \"container\" which permits zinc wrap-around to occur when the jet nozzles are vertically offset.\nEdge build-up of the molten coating metal is eliminated by the practice of the present invention by avoiding oxidation.\nThe zinc coated product produced by the method of the present invention (including spangle minimization) has such excellent surface qualities after temper rolling that it is suitable for use in exposed automotive body panels, appliance applications and the like."}
{"text": "The invention relates to a multilayer film or sheet structure containing at least one barrier layer, at least one predominantly propylene-based layer and at least one predominantly ethylene-based layer, and at least one predominantly ethylene-based adhesive layer.\nPackaging films and sheets are increasingly manufactured by coextrusion systems of growing sophistication. Whereas many years ago, it was common to have machines capable of coextruding three to five layers, such systems today can routinely coextrude nine layers or more, using an equal number of extruders to feed these layers. The new systems allow many types of polymers to be used for better tailoring of multiple properties required in packaging structures, such as barrier, toughness, sealability, gloss, transparency, impact resistance, while lowering overall raw material cost by using less of the expensive components. The polymers themselves many be further blended with other polymers. Additives may be added to the polymers, such as slips or antiblocks, antifog, anti-oxidants, color masterbatches, and the like. Further, the increasing complexity of the material combinations also requires more complicated logistics and handling to ensure that the right materials are delivered to the right extruders during production and to maintain proper inventory economically.\nMany packaging films and laminates contain layers of polar polymers such as polyamide and ethylene vinyl alcohol to provide gas and flavor barriers. With the increasing number of extruders in coextrusion systems it is not uncommon to not only have an ethylene vinyl alcohol layer or a polyamide layer, but also combinations of these layers. Tri-layer combinations of barrier cores or barrier sandwiches comprising a coextrusion of polyamide-ethylene vinyl alcohol-polyamide are widely employed as a method to improve formability in films and sheets. In many cases, the sealant side of the coextrusion construction, that is, the layers between the barrier core and the packaged product are frequently ethylene polymer-based to make use of their good sealing properties. Examples of such sealants are ethylene ionomers, ethylene acid copolymers, ethylene vinyl acetate copolymers, ethylene α-olefin copolymers. In many cases, the exterior side of the coextrusion construction, that is, the layers between the barrier core and the outer layer of the package that does not contact the packaging contents, is composed of propylene polymers. The propylene polymers provide stiffness, moisture barrier, and they also provide a higher temperature resistance that prevents the film laminate from being deformed when contacted against a sealing bar during heat sealing of such laminates. The temperatures of the seal bars can be very high, particularly for thicker laminates, as heat must be transferred from the seal bar to the internal sealant layer within the short contact times necessary for profitable commercial production. For that reason, it can also be common for polyamide to be employed as the outermost exterior layer of the film laminate with polypropylene layers internal to that layer on the exterior side of the coextrusion construction. Consequently, as many as three or more tie layers may be employed in nine or higher layer coextrusions. It would be ideal if a single adhesive resin could be used for all adhesive layers to avoid mistakes with charging the wrong adhesive into the extruder hoppers, and to minimize inventory complications with maintaining multiple resins.\nGenerally speaking, high crystallinity polypropylene (PP) and polyethylene (PE) polymers are not compatible with each other. Most crystalline polyethylene resins do not adhere well to polypropylene and vice-versa. Ethylene copolymer-based adhesives may be employed, such as ethylene vinyl acetate and ethylene methyl acrylate adhesives, but these adhesives do not have sufficient high temperature resistance and may fail when packages are subject to cooking conditions, particularly those at high temperatures and long times. Lower density ethylene-α-olefin copolymer or higher comonomer propylene-α-olefin copolymer based adhesives may also be used, but they suffer from the same temperature issues as they are also lower crystallinity resins because of their high comonomer. They are also more expensive than their lower or no-comonomer polyethylene and polypropylene resin counterparts.\nU.S. Pat. No. 4,405,667 discloses a retortable packaging pouch where the sealant layer is polyethylene and the outer layer is a polyamide. Internal layers of polyamide and ethylene vinyl alcohol are described. The pouch can contain up to 9 layers where one or more of the layers is polypropylene-based. The sealant PE is bonded to the propylene layers by a blend of linear low density polyethylene and a propylene ethylene copolymer. The layer directly adhering to the nylon is an anhydride-modified polypropylene.\nU.S. Pat. No. 4,735,855 discloses a thermoformable film laminate in which the sealant layer is an ionomer or ethylene copolymer and is directly adhered to both a nylon, ethylene vinyl alcohol, nylon barrier core and a polypropylene outer layer. The tie layers bonding to the barrier, the sealant and the outer polypropylene layers may be identical or different from each other and may include a wide range of anhydride grafted polyolefins. The examples demonstrate only anhydride modified ethylene vinyl acetate polymers as adhering the sealant and polypropylene to the barrier polymer.\nU.S. Pat. No. 5,643,999 discloses a coextrudable adhesive composition, for bonding polyolefin film to ionomer film, produced by blending polyethylene with a very low density polyethylene (VLDPE) together with a hydrocarbon elastomer. However, such adhesive composition does not adequately bond polyolefin film layer to barrier layer, for water, oxygen or other substances, such as polyamide or ethylene vinyl alcohol (EVOH) layer.\nU.S. Pat. No. 6,184,298 discloses an adhesive composition of an anhydride modified polyethylene, an unmodified polyethylene and a styrene-based elastomer as being useful in multilayer structures containing nylon or EVOH as the barrier, and having polypropylene as one substrate and polyethylene as another.\nU.S. Pat. No. 7,713,636 discloses a multilayer film comprising a PP core layer, an outer barrier skin layer and an optional second skin layer on the opposite side of the core layer than the barrier skin layer, where the optional second skin layer is selected from the group consisting of a PE polymer or copolymer, a PP polymer or copolymer, an ethylene propylene copolymer and an ethylene vinyl acetate polymer. A tie layer is disposed between the skin and core layers, and comprises at least 10 w % grafted propylene-based polymer containing 5 to 24 weight % ethylene and having a DSC melting temperature of 120° C. or less and a heat of fusion of 75 J/g or less.\nMore recent commercial advances in α-olefin copolymerizations using single site or metallocene catalysts in C2 to C8 α-olefins have resulted in copolymers that have a more blocky characteristic, where the block segments themselves are ethylene α-olefin copolymers and different blocks are distinguished by different comonomer content to tailor the block crystallinity. Such copolymers have been found to improve adhesion between polypropylene and high density polyethylene when used as the tie layer material in experiments conducted by Case Western University. (P Dias et al, Polymer 2008; 49:2937-2946 and Y. Lin et al, Polymer 2011; 52:1635-1644). Here, the polymers are used at 100% in the tie layer on a lab line extruding a small tape of about 2 mm thick and 12 mm wide. Line speeds were not disclosed.\nIt is particularly desired that an adhesive composition that provides high bond strength to multilayer barrier films where one tie layer type can be used to bond polypropylene polymers and copolymers to a polar barrier substrate as well as to polyethylene polymers, ethylene α-olefin copolymers and ethylene copolymers such as ethylene acid ionomers, ethylene acid copolymers, ethylene vinyl acetate and ethylene acrylate copolymers. In particular, it is challenging to find a common tie layer to bond polypropylene and ethylene acid ionomers of sodium and highly neutralized ethylene acid ionomers within a multilayer barrier structure."}
{"text": "This invention relates generally to tools and relates, more particularly, to means for use when pounding, or driving, a post into the ground.\nIn my earlier U.S. Pat. No. 6,719,261, I described an apparatus for supporting a post, such as a metal t-post, which is to be driven into the ground by an individual. Briefly, that apparatus includes a platen base upon which an individual can stand while the post is being driven into the ground, and the base defines a through-opening which is disposed between the feet of the individual when the individual stands upon the base. A cavity-providing guide is mounted upon the base for slidably receiving a post inserted endwise therein and so that the cavity of the guide means is positioned in registry with the through-opening of the base. When an individual stands upon the base with his feet disposed on opposite sides of the through-opening and a post is positioned within the guide means, the post is supported by the guide means in a steady condition for being driven into the ground.\nA limitation associated with the apparatus of my earlier patent relates to the fact that once the platen base of the apparatus is positioned upon the ground, the orientation of the post which is desired to be driven into the ground relative to the surface of the ground cannot be appreciably altered. In other words, in order to make a significant adjustment in the orientation of a post desired to be driven into the ground once the platen base of my patented apparatus is placed upon the ground, the position of the base will have to be re-adjusted.\nAccordingly, it is an object of the present invention to provide a new and improved post support apparatus which is capable of supporting a post while it is being driven into the ground and which accommodates a substantial adjustment in the orientation of a post desired to be driven into the ground without requiring that the base of the apparatus be moved.\nAnother object of the present invention is to provide such an apparatus which is uncomplicated in structure, yet effective in operation."}
{"text": "1. Field of the Invention\nThe invention relates to a process for bending glass sheets and a device for supporting a glass sheet heated to the forming temperature in a horizontal position, said device consisting of an enclosure in which a partial vacuum can be established and of a contact surface provided with suction orifices directed downward and against which the glass sheet is held by a suction effect.\n2. Discussion of the Background\nVarious processes for producing bent glass sheets are known, in particular automobile glazings, which use a support device consisting of an enclosure in which a partial vacuum is established and of a contact surface provided with suction orifices directed downward. The support device is used to lift the glass sheets, heated to the bending temperature in a tunnel furnace, above the horizontal conveyor belt and deposit them on a bending tool or on another conveying device. Various embodiments of these support devices, for example, have been proposed in the patent publication Nos. DE-2 000 271, EP-3391, EP-3392 and EP-182 638. In all known embodiments, the contact surface--against which the glass sheet is held by a suction effect--is a plate, of flat surface or slightly bent, provided with vertical bores or ducts leaving the contact surface and ending in an enclosure under a partial vacuum that can be regulated. As a general rule, the partial vacuum is obtained by means of a fan placed directly on this support device. The support device is mobile along the vertical plane and optionally in a horizontal plane.\nThe plate provided with suction orifices that the support devices according to the art comprise is generally of a specific ceramic. The ceramic plates are rather expensive, because all the ducts must be provided.\nWhen these support devices are used, the entire top face of the glass sheet comes in contact with the lower face of the suction plate. The feel of the top face of the glass sheet is always unfavorable if the glass sheets are sensitive to direct contact. When the top face is provided with a sensitive overlayer, for example, with a layer of printed enamel which is then in the molten state considering the temperature of the glass, these support devices cannot be used because the layer of molten enamel sticks to the contact surface and is then destroyed or at least damaged by the mechanical contact.\nIt is also known how to shape a glass sheet conducted by a roller conveyor under an upper bending form and lifted by an ascending current of hot air which flattens the glass sheet against the upper bending form. The glass sheet is then retrieved by a tempering frame which transfers it between blowing boxes. This process is described in a detailed manner in French Patent No. 2 567 508. The upper bending form, in this case, does not need to be pierced and is therefore less expensive. Moreover, the current of hot air acts on the entire lower face of the glass, therefore, in a way there is a pneumatic pressing which assures a very good shaping quality. However, in this process the energy expenditures are rather great because the hot air current is established for a long period of time."}
{"text": "It has long been known that a plurality of nanocrystallites in silicon carbide (SiC) would give rise to an enlargement of the energy gap of the SiC shifting any emitted light towards the UV region because of quantum confinement, which allows the relaxation of momentum selection rules by confining the charge carriers spatially, thus allowing direct band gap transitions. See U.S. Pat. No. 5,376,241 entitled, “Fabricating Porous Silicon Carbide” by A. D. Kurtz et al., which issued on Dec. 27, 1994 and is assigned to Kulite Semiconductor Products, Inc, the assignee herein. In that particular patent, it teaches the formation of porous SiC, which is formed under electrochemical anodization. The patent also describes the production of the semiconductor through the use of ultraviolet light to illuminate the surface of the semiconductor. In this manner, by controlling the light intensity, the potential and THE doping level, a porous layer is formed in the semiconductor and thus one produces porous SiC. The porous SiC can be employed for UV light sources, such as LEDs and diode lasers. Porous SiC can be utilized as a filtering chemical process to be used to provide heterojunction devices.\nSee U.S. Pat. No. 5,376,818 entitled, “Large Area P-N Junction Devices Formed from Porous Silicon” issued on Dec. 27, 1994 to A. D. Kurtz et al. and assigned to the assignee herein. That patent shows the formation of porous SiC, which is produced under electrochemical anodization. The patent teaches that when a potential is applied to the semiconductor and ultra-violet light illuminates the surface of semiconductor, one can control the light intensity, the potential in doping level, to form a microporous structure in the semiconductor to thus produce porous SiC. The structure enhances the quantum confinement of energetic carriers and which device is highly sensitive to stress.\nReference is also made to U.S. Pat. No. 5,834,378 entitled, “Passivation of Porous Semiconductors for Improved Optoelectronic Device Performance and Fabrication of Light-Emitting Diode Bases on Same”. The patent issued on Nov. 10, 1998 to A. D. Kurtz et al. and is assigned to the assignee herein. That patent describes a method which improves the photoluminescent performance of a porous semiconductor. According to the method, a monolayer of passivating material is generated on a pore wall of the porous semiconductor to passivate the porous semiconductor. This layer substantially eliminates dangling bonds and surface states which are associated with the porous semiconductor layer. The resulting passive porous semiconductor layer exhibits a quantum efficiency of approximately five percent. It is indicated that one monolayer of passivating material can be an oxide which is generated by placing the bulk semiconductor substrate into a furnace. Also described is a heterojunction light emitting device which employs a passivated porous semiconductor layer.\nU.S. Pat. No. 5,939,732 issued on Aug. 17, 1999 and is entitled, “Vertical Cavity Emitting Porous Silicon Carbide Light Emitting Diode Device and Preparation Thereof” and is assigned to the assignee herein and is invented by A. D. Kurtz et al. That patent teaches a multi-layered light emitting device, which has an active light emitting layer of porous silicon carbide and a sequence of layers of porous SiC underneath which serve as a quarter wavelength multi layer mirror. In this manner, one obtains electroluminescent emission of narrow visible light in the deep blue to UV range in a highly directed pattern. Thus, as indicated above, the nanocrystallites in SiC give rise to an enlargement of the energy gap and shifts emitted light towards the UV region. The same effect has also been demonstrated in silicon. Moreover, when LEDs are made from such materials, the emitted light is shifted towards the UV, the shifting inversely proportional to the size of the nanostructure. It is also well-known that the width of the energy gap may also be effected by the application of stresses (see for instance, deformation potentials). The use of deformation potentials as effecting the energy gap is well-known and is text book material. Thus, it is indicated and known that the effect of stress can cause a change in the frequency of emitted light of an LED or the light resonance of the structure.\nIn graphite, a similar effect can occur. Normal graphite is a semi-metal, but in a nanostructure it can be a conductor or a semiconductor. For example, see an article entitled, “Nanotubes for Electronics”. This article appeared in Scientific American in the December 2000 issue. The article describes nanotubes and is written by Phillip G. Collins and Phaedon Avouris. In that article it is clear that nanotubes are utilized because of their unique electronic properties. Carbon nanotubes essentially can be used to perform the same function as silicon does in electronic circuits, but at a molecular scale, where silicon and other semiconductors do not work. In particular, when the dimensions of the nanotube are of the same order of magnitude as the electron wavelength, then these quantum effects can occur at those dimensions.\nSee also an article entitled, “Cavity Quantum Electrodynamics” which appeared in Scientific American in April, 1993 by Serge Haroche and Jean-Michel Raimond. This article explains the operation of atoms and photons and their behavior in small cavities. The article shows that new sensors can be developed utilizing such techniques.\nIn any event, because of the function and operation of nanotubes, it has been determined that application of stress can change a conductor to a semiconductor by changing the energy gap where the quantum confinement leads to a large change in the electrical properties. Essentially, the electrical properties of nanostructures, such as nanotubes, which exhibit quantum confinement can be changed by the application of various stresses, thus leading to a means of measuring such stresses.\nTherefore, the present invention contemplates the formation of a high frequency, high temperature sensor which employs nanotubes and where the nanotubes are subjected to pressure or other external conditions to vary the electrical properties of the nanotubes according to a desired application.\nThe nanotubes, which are conductive, will respond to an applied pressure to produce an output voltage or current proportional to the amount of pressure applied. These devices, because of their high temperatures properties, as well as because of the electrical properties provide improved transducer structures or sensor structures useful in generating sensors in the optical and microwave ranges."}
{"text": "I. Field of the Invention\nThe present invention generally relates to material structures that include ductile, lightweight, high fatigue strength and high fracture strength characteristics. More particularly, this invention relates to a system and method for providing ease of mobility and protecting a host against an impacting ballistic element.\nII. Detailed Description of the Prior Art\nMany typical examples of systems, such as airplanes, spacecraft, cars, tanks, boats, buildings and personnel are subjected to ballistic impact by projectiles. For example, among others, a ballistic element includes a bullet directed at personnel, an artillery shell impacting a tank, a missile fired at a ship, a high speed particle hitting a satellite, and flying debris impacting an automobile or window. Impacts by ballistic elements often result in catastrophic or terminal damage to these systems. To mitigate or prevent impact damage, structural systems are often armored with laminates, i.e. protective layers of material deposited on the system, or are structurally composed of an anti-ballistic material.\nUnfortunately, armored systems are often heavy and slow moving in that the mass density of armor typically increases with increasing ballistic protection. Thus, there exists an increased risk of detrimentally compromising the overall performance of a system as more armor is added to the system. For example, to adequately protect military aircraft from potential ballistic threats, weight from added armor protection could prevent an aircraft from maneuvering, mobilizing over long distances, or even limit the performance of basic flight operations such as taking off and landing for example. In another example, to protect a satellite system from impacts from celestial projectile particles, satellite components are often plated with armor but any increased weight from the armor would dramatically increase the cost for delivering that satellite payload into space.\nAll of these existing types of armor suffer from the inability to provide a lightweight material having high fatigue and fracture strength. Many typical examples that provide protection against ballistic elements include either a thick and heavy armor-type or a lighter armor type that provides easy mobility but with less ability to sustain a ballistic impact than heavy armor. Heavy armor types typically provide greater fatigue and fracture strength whereas lightweight composites provide greater mobility while often compromising fracture and fatigue strength.\nIllustratively, tanks, ships, and armored personnel vehicles typically include armor often having metallic or other high density components that render such devices as bulky, heavy, and slow moving. However, to provide greater requisite mobility, lightweight armor often includes composite materials that are brittle and sacrifice the fatigue and fracture strength provided by heavy armor. Some examples of lightweight armor include protective material for aerospace applications and for personnel such as bulletproof vests and helmets.\nFor example, a typical tank includes heavy, bulky armor for sustaining engagement with a wide range of ballistic elements in combat. However, the weight associated with the armor often makes the tank difficult to move. For example, a tank may encounter difficulty during military maneuvers or as the tank is mobilized to the theatre front. For one specific example, a US-made M1 Abrahams tank, although ideal for sustaining projectiles at high speeds, incorporates a uranium-based armor. During mobilization to the theatre front, this armor prohibits the M1 tank from transport through ideal avenues of air and truck transport exclusively to the slower options of by ship or rail due to the extreme weight of the uranium armored tank. On the other hand, lightweight armor provides for ease of movement but often compromises the ability to sustain a greater range of projectile sizes and/or projectile speeds as that of heavier armor. Armored aircraft, such as for example tactical or strategic aircraft, require lightweight armor for optimal maneuverability through the air. For tactical aircraft, graphite-epoxy composites provide lightweight armored protection against a narrow range of projectile sizes and/or speeds.\nIn short, existing types of armor do not provide for the aspects of lightweight as well as high fatigue and fracture strength. Many other problems and disadvantages of the prior art will become apparent to one skilled in the art after comparing such prior art with the present invention as described herein."}
{"text": "The present invention pertains to an improved irrigator aspirator tip component of the type inserted into the lens capsule of an eye, such as for removing cortical material, washing, cleaning and/or polishing."}
{"text": "An internal combustion engine may be equipped with an air charging system designed for supplying air into the combustion chambers. The air charging system generally includes an intake pipe that conveys air from the ambient environment to an intake manifold in fluid communication with the combustion chambers. An intake valve is generally disposed in the intake pipe. The intake valve has an actuator arranged to move a valve member that regulates the mass flow rate of the air that flows towards the intake manifold.\nThe air charging system may also comprise a turbocharger provided for increasing the pressure of the air in the intake manifold. The turbocharger comprises a compressor located in the intake pipe and a turbine that rotates the compressor by receiving exhaust gasses from an exhaust manifold in fluid communication with the combustion chambers. To regulate the rotational speed of the compressor, the turbine may be a variable geometry turbine (VGT) having an actuator arranged to alter the angular position of the turbine inlet vanes. The air charging system may comprise an additional air compressor located in the intake pipe and driven by an electric motor.\nThe air charging system may further comprise one or more exhaust gas recirculation (EGR) pipes for recirculating a portion of the exhaust gasses from the exhaust manifold back into the intake manifold. Each EGR pipe is generally provided with an EGR valve having an actuator arranged to move a valve member that regulates the mass flow rate of the recirculated exhaust gasses.\nDuring the operation of the engine, the actuators of the air charging system, such as the EGR valve actuator(s), the VGT actuator, the intake valve actuator and the electric motor of the air compressor, are used to regulate a number of output parameters of the air charging system, for example the pressure inside the intake manifold, the oxygen concentration inside the intake manifold and the pressure inside the exhaust manifold, according to performance and emissions requirements.\nTo perform this function, these actuators are conventionally controlled according to separated and uncoordinated control strategies, which enable each individual actuator to be operated in any desired or required position, independently from one another.\nAccordingly, it is desirable to provide improved control strategies for an air charging system. In addition, it is desirable to coordinate the control strategies of the actuators in the air charging system. Furthermore, other desirable features and characteristics of the present disclosure will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "A PWM converter is known as a device which performs switching control performed by PWM to thereby suppress a harmonic current included in a load current for an AC power source driving a motor driving device in order to improve power.\nA PWM converter outputs a square wave AC voltage including high frequencies of several kHz or more on a path between the AC power supply and the PWM converter to perform switching control by the PWM. As the square wave AC voltage includes high frequency components not included in the power supply frequency, this type of PWM converter is generally configured to provide a low pass filter (LC filter) which allows a high frequency square wave AC voltage to pass between the AC power source and the converter in order to remove the high frequency waves.\nFIG. 1 illustrates an example of a basic configuration of an LC filter. FIG. 1 illustrates a configuration where the LC filter 1001 is positioned between the PWM converter 1002 and a three-phase AC power supply 1000.\nThe LC filter 1001 for the PWM converter is configured to be equipped with a so called T-filter which comprises capacitors C1 to C3 whose ends have an inductance of La1 to La3 and Lb1 to Lb3 respectively. Generally two AC reactors 1010 and 1020 for realizing two inductances have similar configurations (the cores are positioned in the same direction). Further, in general, AC reactors are of a natural air cooling type (for example Japanese unexamined patent publication No. 2007-221858 (JP 2007-221858 A)).\nRegarding the inductances L of the two AC reactors 1010 and 1020, in order to suppress the peak of the high frequency current flowing in from the PWM converter 1002 side, the L of the PWM converter side increases and the L of the AC power source 1000 side decreases. Therefore, in general, the volume of the AC reactor B (1020) on the PWM convertor side is comparatively larger than that of the AC reactor A (1010) of the AC power source side. Further, in general, the outer shape of the AC reactor is a parallelepiped about which the coil is wound.\nFurther, a large part of the loss generated in inductance is on the PWM convertor side where high frequency current flows. Therefore, high frequency current only flows through the AC reactor B on the convertor side (1020) which further increases the temperature rise.\nConventionally, due to AC reactors being large and heavy, two AC reactors and a capacitor, etc., constituting an LC filter were not housed in the same casing.\nIf these two AC reactors were to be housed in the same casing, there would be a problem that a useless space would be created due to the difference in the height between the small AC reactor on the power source side and the large AC reactor on the convertor side and the volume of the casing would be large.\nThe object of the present invention is to enable efficient cooling of the AC reactors and provide an LC filter which allows miniaturization and reduction in weight of a casing which houses the LC filter part."}
{"text": "The grinding of wooden frames causes problems since the grain fiber directions of the longitudinal and transverse wooden parts of the frame are different, and it is desirable that the upper surface of a frame wooden part not be ground transversely to its fiber direction.\nTo avoid this disadvantage, band grinding machines are already known in which for a cross-grinding, a transverse band assembly and a broad band assembly are arranged behind one another in the advancement direction. The transverse band first grinds the top of the entire frame. The broad band assembly has two pressure shoes or pressure bars arranged parallel to the advancement direction which are adjustable transversely to the advancement direction so that the broad band can targetly be pressed against the longitudinal wooden parts of the frames to be ground. Therefore, the grinding tracks produced on the wooden longitudinal parts by the transverse band can be removed. A disadvantage of this solution is, however, that the grinding band of the broad band assembly is worn only locally.\nIn a further known cross-grind grinding machine, the broad band by means of an electronically controlled multiple member pressure bar is pushed downwardly in areas corresponding to the longitudinal wooden pieces. Also in this case, there arises the disadvantage that the broad band is only locally worn. Moreover, problems arise in the adjustment of the grinding width to the actual longitudinal wooden piece width, since the individual pressure elements of the pressure bar cannot be made as small as desired and, therefore, the danger exists that the longitudinal wood pieces are either not ground up to their longitudinal edges or that the longitudinal edges are ground round.\nThe invention has as its object the provision of a band grinding machine of the aforementioned kind which permits a desired grinding of the work piece parts parallel to their grain fiber directions with optimal utilization of the used grinding means."}
{"text": "The present invention relates to variable-strength multipole beamline magnets, and more specifically, to a beamline magnet that permits the adjustment of not only the field strength but also the magnetic centerline.\nA number of techniques are available for producing variable-strength magnets. They are especially useful for bending, focusing, and higher-order control of beams in charged particle accelerators. Most charged particle beam accelerators use magnets to control the beam. This is especially true for high-energy accelerators, i.e., relativistic particle accelerators. The magnets affect the beam in ways that are mathematically similar, but not identical, to how optical lenses and mirrors affect an optical beam. In the present description, devices based on pseudo-optical properties of magnets are called beamline magnets.\nCommon beamline magnets are dipoles, quadrupoles, and sextupoles. Dipoles change the direction of the beam as well as provide some focusing or defocusing, like a light pipe with lenses. Quadrupoles focus the beam like a lens. Sextupoles can be used to correct certain types of aberrations. More generally, a beamline magnet with a plurality of poles, including dipoles, quadrupoles, and sextupoles, is termed a multipole magnet. For example, an octupole that uses eight poles is also a multipole magnet, which is suitable for correcting higher-order distortions of the beam.\nMany beamline magnets are electromagnets. In these devices ordinary or superconducting coils are wound around specially shaped poles to generate the desired magnetic field. Adjusting the current passing through the coil(s) controls the magnetic field strength. This has the desirable property that the pole shape controls the field quality. The coils simply supply the magnetomotive force needed to generate the field. Room temperature coils usually need cooling to dissipate the heat generated by the finite resistance of the coils. This is accomplished by using fans, cooling channels, or liquid-cooled copper tubing for forming the coils. When copper tubing is used to form the coils, deionized water is circulated within the tubing while the current flows through the copper. There are a number of limitations to electromagnets. One is that expensive electrical power and additional plumbing are needed to operate these magnets. In addition, an electromagnet has a size limitation because the current densities, with which the power dissipation scales quadratically, are inversely proportional to the magnets\"\" linear dimension. Thus, smaller electromagnets need to use reduced currents to avoid cooling problems, and cannot have strong fields.\nA second, less common type of beamline magnet is made by arrangements of specially shaped magnets. These devices use special arrangements of magnets without poles to produce the desired fields. Sample magnets of this type can be found in U.S. Pat. No. 4,355,236 to Holsinger and U.S. Pat. Nos. 4,429,229 and 4,538,130 to Gluckstern. In these devices, the magnetic field strength is adjusted by rotating rings or disks of magnets. Because of the absence of poles, the magnetic fields of the individual magnets superimpose on each other, which makes analysis of their performance much easier. These magnets also have the advantage that they do not require power supplies to generate currents in the coils or plumbing for cooling the coils as in the electromagnets. However, the field quality produced by these magnets is inferior to that produced by electromagnets. Any mechanical imperfection of the magnets or magnetization nonuniformity degrades the magnetic field quality.\nA third type of beamline magnet uses poles to produce a high-quality field like the one produced by an electromagnet, but uses permanent magnets in place of the coils used in an electromagnet. A sample device of this type can be found in U.S. Pat. No. 4,549,155 to Halbach, wherein the field strength is adjusted by rotating magnets. The rotation of magnets, however, causes the field strength to vary nonlinearly and sinusoidally as a function of a rotating angle, which makes it difficult to adjust the field strength with high precision. Another example of the type of beamline magnet using poles and permanent magnets can be found in U.S. Pat. No. 2,883,569 to Kaiser et al. In this patent, a flux shunt selectively slides over a portion of a cylindrical magnet to short out a varying amount of the magnetic field. This design, though, is intrinsically less efficient because there is a major magnetic flux leakage path between pairs of poles. In addition, this design also produces a nonlinear field adjustment, which is not desirable for high-precision strength adjustment. Yet another example of this type of beamline magnet uses cylindrical magnets that are individually rotated about their axes of symmetry. For these designs, there is one rotating magnet for each pole. The field strength is varied by adjusting the angular position of each magnet with respect to each pole. As before, this style of magnet produces a sinusoidal variation in the magnetic field strength and it is difficult to remove backlash in the rotational system to achieve precise adjustment of the field strength. In addition, many applications require a field strength setting (xcex94B/B) of {fraction (1/10000)} (0.01%). This implies extremely fine angular resolution: the angular encoders need to have resolutions of {fraction (1/50000)} radians, or approximately 300,000 encoder ticks in 360 degrees, which would be extremely difficult to obtain, if not impossible.\nA need exists for a beamline magnet which does not require power supplies or plumbing, and yet produces a high-quality field. Preferably, such a beamline magnet is capable of achieving nonsinusoidal field strength adjustment to allow for high precision adjustment.\nThe present invention provides a multipole beamline magnet that is capable of selectively adjusting magnetic field strength and a magnetic centerline. Specifically, the beamline magnet includes a plurality of stationary poles formed of ferromagnetic material and one or more permanent magnets that are disposed between the plurality of stationary poles. Each of the permanent magnets supplies magnetomotive force to two adjacent stationary poles, so that the poles produce a magnetic field in a central space defined by the poles. A mechanical axis of the beamline magnet extends through the central space perpendicularly to the plane defined by the magnets and the poles. The beamline magnet further includes a linear drive for moving the permanent magnet(s) along radial lines perpendicularly to the mechanical axis, i.e., radially inward or outward with respect to the mechanical axis. Thus constructed, the beamline magnet produces a high-quality field using its stationary poles, and further allows for precise adjustment of the magnetic field strength and the magnetic centerline by collectively or selectively moving the permanent magnets.\nIn accordance with one aspect of the invention, the beamline magnet further includes a pair of nonmagnetic end caps that are provided to sandwich the poles and the magnets. In one embodiment, at least one of the end caps defines one or more guide channels for movably mounting the one or more permanent magnets, respectively. The guide channels are provided for greater control of the linear movement of the magnets.\nIn accordance with another aspect of the invention, the beamline magnet further includes a pair of ferromagnetic shield plates mounted on the nonmagnetic end caps, to thereby sandwich the nonmagnetic end caps, which in turn sandwich the poles and the magnets. The shield plates are used to effectively eliminate magnetic interactions between the beamline magnet and nearby instruments or other beamline magnets.\nIn accordance with yet another aspect of the invention, the beamline magnet further includes a magnetic field sensor arranged to determine the strength of the magnetic field in the central space defined by the stationary poles. The sensed magnetic field strength data may then be used to control the linear drive for selectively or collectively moving the permanent magnets.\nIn accordance with still another aspect of the invention, the beamline magnet further includes a beam position sensor arranged to sense the location of a charged particle beam in the central space defined by the stationary poles. The sensed beam position may then be used to control the linear drive for selectively or collectively moving the permanent magnets to adjust the magnetic field strength or magnetic centerline.\nIn accordance with still another aspect of the invention, the beamline magnet includes a means of passive temperature compensation for maintaining the magnetic field strength substantially constant regardless of any changes in the operating temperature. Specifically, ferromagnetic materials having a low Curie temperature are magnetically coupled to the permanent magnets in a parallel flux shunting configuration to compensate for temperature-dependent flux variation of the permanent magnets. At a low temperature, the permanent magnets are stronger than at a high temperature, and thus could supply more flux in the central space than at a high temperature. At a low temperature, though, the ferromagnetic materials shunt a larger fraction of the available flux away from the central space than they do at a high temperature. Consequently, the resulting flux in the central space is substantially the same at both low and high temperatures; at a low temperature, the magnets are stronger but more flux is shunted away from the central space, and at a high temperature, the magnets are weaker but less flux is shunted away from the central space. With proper choice of the ferromagnetic material, its dimensions and location, the magnetic field strength can be maintained at an essentially constant level despite changes in the operating temperature.\nIn accordance with still another aspect of the invention, the beamline magnet includes a means of passive temperature compensation to correct for thermally induced shifts of the magnetic centerline. Centerline shifts can be caused by various thermal reasons, for example, by thermal expansion or contraction of all the materials in the beamline magnet, temperature dependence of the magnetic properties of the permanent magnets, and temperature induced movement of a support platform on which the beamline magnet is mounted. According to the present invention, thermal compensation of centerline shift is achieved by coupling different amounts of temperature compensating material (i.e., ferromagnetic material having a low Curie temperature) on each magnet. With proper choice of the material, its dimensions and location, the magnetic centerline can be maintained at an essentially constant location despite changes in the operating temperature.\nIn accordance with still another aspect of the invention, the beamline magnet further includes electromagnetic corrector coils to make small adjustments to the magnetic centerline and/or the magnetic field strength. One or more corrector coils are strategically placed to selectively supply a predetermined amount and polarity of magnetomotive force to one or more stationary poles. Adjustment using the electromagnetic corrector coils is achieved by merely modifying wiring of, and the current passing through, the coils, and hence the adjustment is quick and precise. For fine-tuning the field strength and/or the magnetic centerline, electromagnetic adjustment may be more advantageous than the mechanical adjustment of the present invention using the linear movement of the permanent magnets.\nIn accordance with still another aspect of the invention, the beamline magnet includes a plurality of poles and a plurality of permanent magnets. The poles and the magnets may be provided in equal numbers, and may be arranged equiangularly over 360xc2x0. The poles may be made of various materials and in various shapes. All the poles in a beamline magnet may be fabricated the same, or differently from each other. Likewise, the permanent magnets may be made of various materials, in various shapes, and having various magnetization directions. All the permanent magnets in the beamline magnet may be fabricated the same or differently from each other. Furthermore, each of the permanent magnets may be formed of a plurality of submagnet portions having the same or different shapes or properties. The shapes and properties of each pole and each permanent magnet (or submagnet portion) are determined so as to produce the desired magnetic field distribution according to each application.\nIn accordance with still another aspect, the beamline magnet of the present invention further includes one or more stationary auxiliary magnets positioned between the central space defined by the poles and the one or more permanent magnets, respectively. In other words, the auxiliary magnets are arranged radially inward of the permanent magnets with respect to the mechanical axis. The auxiliary magnets remain fixed while the permanent magnets disposed radially outward of the auxiliary magnets are moved.\nIn accordance with a further aspect, the beamline magnet of the present invention includes a ferromagnetic tuning shim. For example, the shim may be attached to the stationary auxiliary magnets, moving permanent magnets, poles, end magnets, or the nonmagnetic end caps. Shims serve to compensate for field errors produced due to imperfection in fabricating the permanent magnets and/or the poles.\nThe present invention further provides a method of selectively adjusting a magnetic field in a multipole beamline magnet. The method includes three steps. First, a plurality of stationary ferromagnetic poles are provided. Second, a plurality of permanent magnets are arranged between the plurality of stationary ferromagnetic poles, so that each of the permanent magnets supplies magnetomotive force to two adjacent stationary ferromagnetic poles. As a result, the stationary ferromagnetic poles produce a magnetic field in a central space defined by the stationary ferromagnetic poles. A mechanical axis of the beamline magnet is defined to extend through the central space, perpendicularly to the plane defined by the magnets and the poles. Finally, the one or more permanent magnets are moved perpendicularly to the mechanical axis.\nThe method may be applied in various ways to achieve the desired adjustment to the magnetic field, such as adjusting the field strength and the magnetic centerline. In a general case, the magnets are individually moved to selectively adjust the magnetic field strength and the magnetic centerline.\nIn a more special case, one may apply the method to adjust the strength of the magnetic field without changing the field distribution. This may be done, for example, by collectively moving all the permanent magnets in a radially inward or outward direction relative to the mechanical axis so as to uniformly increase or decrease the magnetic flux coupling to all the poles. The strength adjustment may be linear, thus allowing for high precision adjustment.\nAs another special case, one may adjust the magnetic centerline without changing the field strength. This may be done, for example, by moving a pair of opposing permanent magnets that are 180xc2x0 apart in one direction. Such movement merely translates (i.e., shifts in parallel) magnetic flux lines, and in effect linearly moves the magnetic centerline.\nThe present invention offers various advantages. First, the beamline magnet of the present invention does not require power supplies or plumbing, and yet produces a high-quality field due to the use of stationary poles. Second, the invention allows for linear adjustment of the field strength and the magnetic centerline, which in turn permits high precision adjustment of the field strength and the centerline. Third, in the present invention the magnets are moved linearly to make various adjustments, as opposed to being rotated, thus the precise adjustment of the magnets is made easier. This permits extremely accurate adjustments of the field strength (0.01%) and the magnetic centerline (microns) with commercially available linear encoders having 1-20 micron resolution. As discussed above, designs that use rotary motion typically require angular resolutions of approximately 300,000 encoder ticks in 360 degrees for 0.01% accuracy. This is not easily achieved with any commercial encoders.\nFourth, the present invention is versatile in permitting various adjustments of the magnetic field. For example, the present invention may be used to adjust the field strength without changing the magnetic centerline, or adjust (shift) the magnetic centerline without changing the field strength. Fifth, the versatile field adjustment capability described above may be readily applied to compensate for any errors in the magnetic properties of the beamline magnet (i.e., magnetic field strength, magnetic centerline, and magnetic field distribution) introduced during fabrication of the beamline magnet. For example, if the permanent magnets have differing strengths, then they can be moved linearly to compensate for the differences. If the magnetization direction of the permanent magnets is nonuniform, then the tuning shims can be used to compensate. Likewise, imperfections in the pole shapes or poles\"\" magnetization properties can be compensated for by combinations of linear motion of the permanent magnets and the use of ferromagnetic tuning shims. Furthermore, when electromagnetic corrector coils are provided, fine adjustments of the field strength or the magnetic centerline can be readily achieved by selectively wiring and passing a current thorough the coils. Thus, the present invention is highly tolerant to variations in the quality of the magnets and/or poles, thereby reducing the overall cost of manufacturing.\nLastly, the construction of the beamline magnet is such that it allows one to access the central space of the beamline magnet by removing one or more permanent magnets. This advantageously permits the beamline magnet to receive an electron beam sensor adjacent the central space for monitoring the behavior of the electron beam passing through the beamline magnet."}
{"text": "The present invention relates to a novel human protein which is involved in the targeting of proteins towards proteasome degradation pathways. This protein, called h-xcex2TrCP, is capable of interacting notably with the Vpu protein of HIV-1 virus and with the cell proteins IxcexaB, xcex2-catenin and Skp1p.\nThe degradation of proteins by proteasome, a multiprotein complex present in all cells, is involved in numerous essential cell phenomena such as the control of cell proliferation, the renewal of proteins and the removal of incorrectly folded proteins, particularly in the endoplasmic reticulum (CIECHANOVER A., Cell, 79, 13-21,1994). Numerous viruses, like HIV-1 virus, which degrades CD4 via one of its proteins Vpu (MRONO D., Cell, 82, 189-1992, 1995), exploit these cell pathways of protein degradation, in which the proteins are targeted towards proteasome by various interactions with other proteins before being degraded. To be targeted towards and degraded by proteasome, the proteins must generally be ubiquitinylated beforehand by ubiquitin-ligase complexes. Furthermore, to be ubiquitinylated, the proteins must often undergo modifications such as phosphorylations (CIECHANOVER A., Embo. J., 17, 7151-7160, 1998).\nSeveral other proteins of the xcex2TrCP type are known at the present time.\nthe xcex2TrCP protein of Xenope, described by Spevak et al. (Mol. Cell. Biol., 13, 4953-4966, 1993);\nthe Slimb protein of drosophila, described by Jiang et al. Nature, vol. 391, Jan. 29, 1998); and\nthe KIAA 0696 protein identified by Ishikawa et al. (DNA Research, 5, 169-176, 1998) during a systematic analysis of sequences expressed in the brain.\nJiang et al. showed that the Slimb protein of drosophila is involved in the stability of the Armadillo protein and the signaling of two metabolic pathways essential for development, namely the Hedgehog and Wingless pathways. They also showed that the Slimb protein has a homology of about 80% with the xcex2TrCP protein of Xenope, none of whose functions was described by Spevak et al. As the xcex2-catenin of Xenope or man, which is the homolog of the Armadillo protein of drosophila, seems to be targeted towards proteasome degradation pathways in the absence of signaling of the Hedgehog and Wingless pathways, said authors suggest that, in man, the genes coding for the homologs of Slimb could be involved in the proteolytic degradation of xcex2-catenin, a protein which acquires oncogenic properties when it is not degraded (POLAKIS P., Biochim. Biophys. Acta, 1332, F127-47, 1997).\nHowever, despite the fact that conservation of the Wingless and Hedgehog pathways in vertebrates is important, it is not certain that the functions of the homologous proteins will be totally conserved. Moreover, there are numerous examples which show that there are always significant differences between species.\nAlso, solely on the basis of genetic studies, Jiang et al,. established the involvement of Slimb in the control of the Wingless and Hedgehog pathways in drosophila. Proof that this control is dependent on a direct interaction between Slimb and Armadillo, for example, has neither been sought nor found.\nThe protein according to the invention, called h-xcex2TrCP, is capable of interacting with virus proteins or cell proteins which can act as mediators or be degraded by proteasome. In particular, the h-xcex2TrCP protein is capable of interacting notably with the Vpu protein of HIV-1 virus and with the cell proteins IxcexaB and xcex2-catenin.\nIt is particularly useful for screening therapeutic agents such as, in particular, antitumoral, antiviral, anti-inflammatory and anti-Alzheimer agents.\nThe Vpu protein is a small membrane protein of 81 amino acids which is expressed by the majority of isolates of HIV-1 virus but not by those of the considerably less pathogenic FUV-2 virus or by those of SIV simian virus (COHEN et al., Nature, 334, 532-534, 1988, and STREBEL et al., Science, 2, 1221-1223, 1988).\nOne of the functions of the Vpu protein is its capacity to induce degradation of the CD4 protein, a cell receptor of HIV-1 virus, so it participates in reducing the expression of the CD4 receptor on the cell surface (Willey et al., J. Virol., 68, 1207-1212, 1994).\nIt is also known that the two phosphorylation serines of the Vpu protein, located in positions 52 and 56, are essential for the degradation of CD4 induced by Vpu (MARGOTIN et al., Virology, 223, 381-386, 1996). Moreover, during the process of infection by HIV-1 in the absence of the Vpu protein, the Gp160 envelope precursor and the newly synthesized CD4 protein combine in the endoplasmic reticulum to block the maturation of the Gp160 protein (BOUR et al., J. Virol., 65, 6387-6396, 1991. Degradation of the CD4 receptor mediated by the Vpu protein is essential for releasing the viral envelope protein which is held in the endoplasmic reticulum by being bound to CD4 through interaction with the Gp120 subunit, and for allowing the normal maturation of the envelope into the plasmic membrane and subsequently its integration into the virus particles, rendering them infectious. Recent studies have demonstrated the fact that degradation of the CD4 receptor mediated by the Vpu protein is sensitive to specific proteasome inhibitors and is dependent on the presence of an xe2x80x9cintact ubiquitinylation machineryxe2x80x9d (FUJITA et al., J. Gen. Virol., 78, 619-625, 1997).\nThus the Vpu protein participates in absolutely critical functions for assuring the production of large numbers of infectious virus particles, since it acts not only on the products of the gag gene, i.e. on the structural proteins, to increase the release of the virus particles, but also on the products of the env gene to allow the maturation of the envelope protein following degradation of the CD4 receptor. In 1996, MARGOTTIN et al. (supra) showed that the interaction between Vpu and CD4 took place via their cytoplasmic domain and that this interaction was not sufficient to trigger degradation of the CD4 receptor.\nThe Skp1p protein is a cell protein involved in the targeting of proteins towards proteasome degradation pathways, which depends on the ubiquitinylation of the proteins (PICKART C. M., The Faseb Journal, 11, 1055-1066,1997).\nBAI et al. (Cell, 86, 263274, 1996) showed that the Skp1p protein was necessary for ubiquitin-mediated proteolysis and that this degradation took place due to the interaction of Skp1p with proteins containing a unit called F-box.\nThe Skp1p protein is an essential factor in the targeting of cell cycle regulatory proteins by proteasome. Targeting of the degradation of these regulators is particularly necessary when the cell cycle enters the S phase of DNA synthesis (PAGANO M., The Faseb Journal, 11, 1068-1075, 1997). Recent studies showed that the Skp1p protein and F-box proteins are the essential elements of high-molecular complexes called SCF (Skp1p-Cullin-F-box-protein complexes). These SCF complexes play the role of enzyme E3; through their ubiquitin-ligase activity, they allow the last step of the ubiquitinylation of substrate proteins, which are thus targeted towards degradation by proteasome (HOYT A., Cell, 91, 149-151, 1997). It is further pointed out that no Skp1p homolog has yet been identified in drosophila.\nThe IxcexaB protein, which exists in different forms (xcex1, xcex2, xcex5), is the major inhibitor of the NFxcexaB transcription factor, keeping it in the form of an inactive complex in the cytoplasm (Beg A. et al., Genes and Dev., 7, 2064-2070, 1993). After stimulation of the cells by factors such as interleukin-1 (IL1) and tumor necrosis factor (TNF), the IxcexaB protein is phosphorylated on serine residues S32 and S36. This phosphorylation leads very rapidly to the ubiquitinylation of the protein and to the targeting thereof towards degradation by proteasome. The active NFxcexaB factor, for example in the form of two subunits P50 and P65, is then released and imported into the nucleus, where it will be able to activate a very large number of genes and cause inflammatory phenomena in particular.\nThe xcex2-catenin protein is a cell protein controling the essential signal transduction pathways such as the Wingless pathways, which are very highly conserved in all vertebrates (MILLER et al., Genes and Dev., 10, 2527-2539, 1996, and POLAKIS P., Biochim. Biophys. Acta, 1332, F 127-47, 1997). xcex2-Catenin accumulates in cancerous cells, either as a result of mutations which prevent phosphorylation on serine residues 33 and 37 (mutated xcex2-catenin proteins), or as a result of mutations of its cofactor, the APC protein, which is necessary for its degradation.\nThe accumulation of xcex2-catenin due to its non-degradation leads to its importation into the nucleus and to the activation of genes controlled by TCF-LEF promoters, causing cell proliferation and transformation phenomena.\nIt was recently shown that mutations of presenilin-1 in patients suffering from Alzheimer\"\"s disease caused a destabilization and enhanced degradation of xcex2-catenin (ZHANG et al., Nature, 395, 699-702, 1998). These authors showed that non-mutated presenilin-1 binds to xcex2-catenin and thereby contributes to its stability. In Alzheimer\"\"s disease, the mutated presenilin is no longer capable of binding to xcex2-catenin, so the latter is degraded more rapidly. The level of xcex2-catenin is considerably reduced in the neuronal cells of patients suffering from Alzheimer\"\"s disease. The loss of xcex2-catenin causes an enhanced apoptosis of the neuronal cells, which would account for the neuronal loss observed in this pathological condition.\nIt is easy to see that there is an urgent need for means of modulating, namely activating or inhibiting, the targeting of proteins towards proteasome.\nA novel human protein involved in the targeting of proteins towards proteasome degradation pathways has now been found which makes it possible to screen modulators of the targeting of proteins towards proteasome.\nThe present invention therefore relates to a novel human protein, called h-xcex2TrCP, which has SEQ ID No. 2 and which is involved in the targeting of proteins towards proteasome degradation pathways.\nThe h-xcex2TrCP protein possesses 569 amino acids and comprises one F-box and seven WD units having the following positions in the sequence SEQ ID No. 2:\nBecause of the homology of this novel protein with the xcex2TrCP of Xenope, a protein containing 0-transducin units and known as xe2x80x9cbeta transducin repeats containing proteinxe2x80x9d, the protein of the invention is called h-xcex2-ATrCP (human xcex2TrCP).\nVia its WD units, the h-xcex2TrCp protein of the invention is capable of interacting with proteins degradable by proteasome, particularly with virus proteins and cell proteins which possess the phosphorylation unit comprising the amino acids Asp-Ser-Gly-Xaa-Xaa-Ser (SEQ ID NO: 9), in which Xaa is any natural amino acid and in which the serine residues are phosphorylated.\nThe phosphorylation of this unit Asp-Ser-Gly-Xaa-Xaa-Ser (SEQ ID NO: 9) is essential to the ubiquitinylation and subsequent degradation of proteins possessing this type of unit. The h-xcex2TrCP protein is only capable of interacting with proteins containing this unit when the two serine residues are phosphorylated, and it cannot interact with proteins containing a phosphorylation unit in which the serine residues are mutated to non-phosphorylatable amino acids. By interacting with the phosphorylated proteins on this unit, the h-xcex2TrCP protein controls their ubiquitinylation and their screening towards degradation by proteasome.\nThe virus protein Vpu and the cell proteins IxcexaB and xcex2-catenin may be mentioned in particular among these proteins.\nIt has also been found that the h-xcex2TrCP protein interacts via its F-box with the Skp1p protein, so it forms part of a novel SCF complex, SCF-h-xcex2TrCP, which selects certain cell or virus proteins for degradation by proteasome. Through its activity of targeting towards proteasome degradation pathways, the h-xcex2TrCP protein according to the invention acts as cell mediator of the Vpu protein in cells infected with HIV-1 virus.\nWithout wishing to exclude other theories, it is thought that, in cells infected with HIV-1 virus, the virus uses, via the Vpu protein, the SCF complex (of which the xcex2TrCP protein forms part) to induce degradation of the CD4 receptor, which will favor the replication of the virus and the release of the infectious virions.\nThe invention further relates to the peptide fragments of the h-xcex2TrCP protein which result from the addition, deletion and/or replacement of one or more amino acids, said peptide fragments having conserved the activity of interacting with proteins degradable by proteasome, particularly with the Vpu protein of HIV-1 virus, with the cell protein IxcexaB or the cell protein xcex2-catenin and/or with the Skp1p protein.\nThe invention relates in particular to the peptide fragments which comprise at least one of the following amino acid sequences of h-xcex2TrCP:\n251-569,\n292-569,\n292-396,\n292-545 and\n1-291.\nVery particularly preferred peptide fragments are those which are partially or totally devoid of the F-box or those which are partially or totally devoid of the WD units.\nOne particularly preferred peptide fragment is the mutant with residues 32-179 deleted, which is hereafter called xcex2TCPxcex94F.\nThe present invention further relates to the nucleic acid sequences, namely the genomic DNA sequences and the cDNA or mRNA sequences, which comprise or consist of a concatenation of nucleotides coding for the h-xcex2TrCP protein or for any one of its peptide fragments as defined above.\nThe invention relates notably to those nucleic acid sequences coding for the h-xcex2TrCP protein and its peptide fragments described above which are represented by:\na) the cDNA sequence SEQ ID No. 1 coding for said h-xcex2TrCP protein and the cDNA sequences of the nucleic acid fragments coding for said peptide fragments;\nb) the DNA sequences which hybridize with the above sequences under strict conditions;\nc) the DNA sequences which, due to the degeneracy of the genetic code, result from the sequences a) and b) above and code for the h-xcex2TrCP protein or its fragments; and\nd) the corresponding mRNA and DNA sequences.\nThe proteins and peptide fragments according to the invention can be obtained by the genetic engineering technique comprising the following steps:\nculture of a microorganism or eukaryotic cells which have been transformed with the aid of a nucleic acid sequence according to the invention; and\nrecovery of the protein or the peptide fragment produced by said microorganism or said eukaryotic cells.\nThis technique is well known to those skilled in the art. Further details on this subject may be obtained by reference to the following work: Recombinant DNA Technology 1, Editors Ales Prokop, Raskesh K. Bajpai; Annals of the New York Academy of Sciences, volume 646, 1991.\nThey can also be prepared by the conventional peptide syntheses well known to those skilled in the art.\nThe nucleic acids according to the invention can be prepared by chemical synthesis and genetic engineering using the techniques well known to those skilled in the art, as described e.g. by SAMBROOK et al. (supra).\nFor example, the cDNA sequences according to the invention can be synthesized by amplifying the mRNAs of human cells by the PCR (Polymerase Chain Reaction) method, as described e.g. by GOBLET et al. (Nucleic Acid Research, 17, 2144, 1989), using, as primers, synthetic oligonucleotides defined from the DNA sequence SEQ ID No. 1.\nThe amplified nucleic acid fragment can then be cloned by the techniques described by AU SUBEL et al. (Current Protocols in Molecular Biology, chapter 3, supra).\nThe invention further relates to transgenic animals which express a transgene for the h-xcex2TrCP protein of the invention, or transgenic animals in which the xcex2TrCP gene has been invalidated.\nThese transgenic animals or animals in which the h-xcex2TrCP protein gene has been invalidated may be used as models for the in vivo study of perturbation of the cell cycle and proliferation by the absence or overexpression of the gene for the h-xcex2TrCP protein or for truncated or mutated forms of this protein, the Skp1p protein, the Vpu protein, the IxcexaB protein or the xcex2-catenin protein.\nThese transgenic animals are obtained by techniques well known to those skilled in the art, such as those described in Manipulating the mouse embryo; a laboratory manual. HOGAN B., BEDDINGTON R., COSTANNNI F. and LACY E., Cold Spring Harbor laboratory press, second edition, 1994.\nThe preferred animals are mammals such as mice or rats.\nThe invention further relates to the prokaryotic microorganisms and eukaryotic cells transformed with the aid of an expression vector containing a DNA sequence according to the invention. This expression vector, which can be e.g. in the form of a plasmid, must contain, in addition to the DNA sequence of the invention, the means necessary for its expression, such as, in particular, a promoter, a transcription terminator, an origin of replication and, preferably, a selection marker. The transformation of microorganisms and eukaryotic cells is a technique well known to those skilled in the art, who will easily be able to determine, as a function of the microorganism to be transformed, the means necessary for the expression of the DNA sequence according to the invention.\nThe preferred microorganism for the purposes of the invention is E. coli, while the yeast used is preferably Saccharomyces cerevisiae. \nCOS, CHO, SF9, Jurkat and other cells, all of which are listed in the ATCC, may be mentioned in particular as examples of eukaryotic cells which are suitable for the purposes of the invention.\nThe invention further relates to the eukaryotic cells cotransformed with expression vectors containing on the one hand the DNA sequence coding for the Vpu protein, for the Skp1p protein, for the IxcexaB protein or for the mutated xcex2-catenin proteins, and on the other hand a sequence coding for the h-xcex2TrCP protein, said expression vectors also containing means useful for their expression, including in the yeast two-hybrid system.\nThe present invention therefore further relates to the anti-HIV-1 antiviral agents consisting of the peptide fragments of the h-xcex2TrCP protein of the invention which have conserved the properties of interaction of the h-xcex2TrCP protein either with the Vpu protein or with the Skp1p protein. These peptide fragments are devoid of the F-box or the WD units, so they are no longer able to interact with the Skp1p protein or, respectively, the Vpu protein.\nOther antiviral agents, antitumoral agents or anti-inflammatory agents which may be mentioned are antibodies directed against the h-xcex2TrCP protein of the invention and its peptide fragments, said antibodies being a further subject of the invention.\nThese antibodies can be monoclonal antibodies obtained by the well-known method of KOHLER and MILSTEIN (Nature, 256, 495-497, 1975) or polyclonal antibodies obtained by the conventional methods of animal immunization (Antibodies, a laboratory manual. E. Harlow and D. Lane. Cold Spring Harbor laboratory press, 1988).\nFinally, antiviral agents, antitumoral agents or anti-inflammatory agents which may be mentioned are antisense oligonucleotides which block the transcription or translation of the h-xcex2TrCP protein of the invention and which hybridize with a nucleic acid sequence as defined above, said oligonucleotides also forming a further subject of the present invention.\nThese antisense oligonucleotides are prepared by techniques well known to those skilled in the art, such as those described by AUSUBEL et al. (Current Protocols in Molecular Biology, Green Publishing Associates and Wiley-Interscience, New York, 1989, published up to 1997).\nThe peptide fragments of h-xcex2TrCP which possess the F-box or which have conserved both the WD units and the F-box can be used as antitumoral or anti-inflammatory agents.\nThe peptide fragments of h-xcex2TrCP which are devoid of the F-box can be used in gene therapy for the treatment of osteoarticular inflammatory diseases or acute inflammatory syndromes which are accompanied by NFxcexaB activation induced by the massive release of TNFxcex1 during these processes.\nAs illustrated in FIG. 7, the expression of h-xcex2TrCPxcex94F is capable of massively inhibiting, by a factor of about 20, the transcription activation induced by TNFxcex1 Therefore, h-xcex2TrCPxcex94F could act as a potent anti-inflammatory agent in any pathological condition associated with an intense inflammatory reaction due to a release of TNFxcex1. For example, several attempts are currently being made to apply gene therapy to rheumatoid polyarteritis by injecting recombinant viruses into the damaged joints. Vectors expressing h-xcex2TrCPxcex94F can be used in these gene therapy experiments on inflammatory syndromes. These vectors may be of several types (retroviruses, adenoviruses; ANDERSON F., Nature, 392, 25-30, 1998). The expression of h-xcex2TrCPxcex94F may be monitored by its effects on the inhibition of NFxcexaB activation by TNF.\nThe present invention further relates to the use of the h-xcex2TrCP protein, or the nucleic acid sequences coding for this protein or for its peptide fragments, for the screening of therapeutic agents which are capable of modulating the interaction of the h-xcex2TrCP protein with proteins degradable by proteasome, and particularly for the screening of:\nanti-HIV-1 antiviral agents capable of inhibiting the interaction between the h-TrCP protein and the Vpu protein and/or inhibiting the interaction between the h-xcex2TrCP protein and the Skp1p protein;\nantitumoral agents capable of perturbing the regulation of the cell cycle or the protein degradation processes in tumoral human cells by modulating (inhibiting or activating) the interaction between the h-xcex2TrCP protein and the Skp1p protein, and by reactivating the interaction between the h-xcex2TrCP protein and the mutated xcex2-catenin proteins in tumoral cells, or between the TrCp protein and the normal xcex2-catenin protein in tumoral cells devoid of the APC protein;\nanti-inflammatory agents capable of perturbing the activation of the NFxcexaB transcription factor by inhibiting the interaction between the h-xcex2TrCP protein and the IxcexaB protein; and\nanti-Alzheimer agents capable of reducing the degree of degradation of xcex2-catenin in neuronal cells by inhibiting the interaction between the h-xcex2TrCP protein and the xcex2-catenin protein.\nIn fact, by perturbing the Vpu/h-xcex2TrCP and/or Skp1p/h-xcex2TrCP interactions, it is possible:\neither to inhibit the replication and production of HIV-1 virus by infected cells;\nor to inhibit the entry of the cell cycle into the S phase and to have an antiproliferative effect.\nBy perturbing the IxcexaB/h-xcex2TrCP and/or Skp1p/h-xcex2TrCP interactions, it is possible to inhibit the degradation of the IxcexaB protein by proteasome and hence to inhibit the activation of the NFxcexaB transcription factor.\nFinally, by activating the mutated xcex2-catenin/h-xcex2TrCP interaction, it is possible to activate the degradation of the xcex2-catenin which has accumulated in tumoral cells. By inhibiting the xcex2-catenin/h-xcex2TrCP interaction in patients suffering from Alzheimer\"\"s disease, it is possible to reduce the apoptosis of neuronal cells.\nThe antiviral agents can be selected either from random peptide banks on the surface of phages (SCOTT J. et al., Science, 249, 386-390, 1990) or by using random synthetic oligonucleotides according to the technique of the SELEX type (TUERK and GOLD, Science, 249, 505-510, 1990). This technique makes it possible to isolate, from a very large pool of oligonucleotides, those which have a high affinity for the protein of interest, namely the h-xcex2TrCP protein in the present case. They are called aptamers. From these aptamers it will be possible, using the screening method below, to select those which inhibit both the Vpu/h-xcex2TrCP and Skp1p/h-xcex2TrCP interactions.\nThe screening method defined above can be carried out e.g. by using the yeast two-hybrid system in which yeast cells co-expressing the h-xcex2TrCP protein according to the invention and one of the proteins Vpu, IxcexaB, xcex2-catenin or Skp1p are cultivated on appropriate selective media in the presence of the test substance; the selective media are the media commonly used in this field and hence are well known to those skilled in the art.\nThe yeast two-hybrid system is described by FIELDS and SONG in Nature, 340, 245-246, 1989, and in patent U.S. Pat. No. 5,667,973. This two-hybrid system is based on detection of the protein-protein interactions by activation of the His or LacZ reporter gene under the control of Gal4 transcription activation domains in the yeast.\nIn this two-hybrid system, a yeast is cotransformed with a two-hybrid vector containing the cDNA of one of the proteins and a vector containing the cDNA of the other protein, each of said vectors containing either a DNA binding domain or a transcription activation domain. The two proteins are then expressed by the yeast in an appropriate culture medium, for example a histidine-free culture medium. The interaction between the two hybrid proteins allows on the one hand activation of the His3 gene and growth of the yeasts on a histidine-free medium, as well as activation of the LacZ gene, which is disclosed by a color reaction specific for xcex2-galactosidase. It is therefore possible to verify the interaction when the yeasts grow on a histidine-free medium and when a color reaction is observed.\nA further possibility is to use the halo test as described by Valtz and Peter (Meth. Enzymol., 283, 350-365, 1997) to detect whether there is any interaction.\nIt is also possible to use variants of the two-hybrid system, such as the three-hybrid system described by TIRODE et al. (J. Biol. Chem., 272, 22995-22999, 1997) or by COLAS et al. (Nature, 380, 548-550, 1996), in which a peptide inhibiting the interaction can be expressed as a third partner to inhibit the interaction of the other two. A random peptide bank can also be used in this way.\nA further possibility is to use the reverse-hybrid system described by VIDAL et al. (Proc. Natl. Acad. Sci., 93, 10315-10320), in which selection is carried out against an interaction and not for an interaction. In this system, as in the conventional two-hybrid system, it is possible to screen banks of small chemical molecules, including those derived from chemical synthesis, in order to bring yeasts cotransformed with two-hybrid or reverse-hybrid vectors carrying fusions with the Vpu protein, the IxcexaB protein, xcex2-catenin, the h-xcex2TrCP protein or the Skp1p protein into contact with these small molecules in the search for an inhibitor of the Vpu/h-xcex2TrCP, Skp1p/h-xcex2TrCP, xcex2-catenin/h-xcex2TrCP or IxcexaB/h-xcex2TrCP interactions.\nThe screening assays for interaction inhibitors may also be carried out using the conjugative two-hybrid system (FROMONT-RACINE et al., Nature Genetics, 16, 277-282, 1997), the membrane two-hybrid system (BRODER Y. C. et al., Curr. Biol., 8, 1121-1124, 1998) and optionally, if phosphorylations can take place in the bacteria, the bacterial two-hybrid system (KARIMOVA et al., Proc. Natl. Acad. Sci., 95, 5752-5756, 1998).\nThis screening can also be effected in vitro by using one of the proteins Vpu, IxcexaB, xcex2-catenin or Skp1p and the h-xcex2TrCP protein, one of the proteins being immobilized on an appropriate support and the other being labeled by any means used in the methods of detecting biological substances, it being possible for this labeling means to be e.g. a radioactive isotope, a luminescent agent, biotin or a specific antibody.\nOne of the proteins will preferably be immobilized in the form of a fusion protein with glutathione S-transferase (GST) on agarose-glutathione beads or in microtiter plates, the GST serving as an agent for coupling said protein with the beads or with the wells of the plates.\nThis can be done particularly using the scintillation proximity assay (SPA) described by BOSWORTH et al. (Nature, 341, 167-168, 1989) and marketed by Amersham. This assay consists in labeling one of the proteins with a radioactive element, for example tritium, and immobilizing the other protein on magnetic beads or agarose-glutathione beads. The inhibitory effect of the test substances on interactions involving the h-xcex2TrCP protein can easily be detected, without separation of the bound or free radioactive species, according to the protocols described by BOSWORTH et al. (supra).\nAnother possible technique is that of surface plasmon resonances described by KARLSSON et al. (J. Immunol. Methods, 145, 229-233, 1991), using Biacore, marketed by Pharmacia, to isolate the inhibitors of interactions involving the h-xcex2TrCP protein according to the invention.\nThe inhibitory activity of the antiviral agents selected in this way may be verified by assays on CD4+T cells or on chimpanzees infected with HIV-1 virus or SIV Cpz.\nThe antitumoral agents and anti-inflammatory agentsxe2x80x94ligands of the h-xcex2TrCP protein of the inventionxe2x80x94can also be isolated by the two-hybrid techniques or related techniques or by interaction in vitro with combinatorial banks of peptides or other chemical products, as described above.\nThe specificity of the antiviral, antitumoral or anti-inflammatory agents selected by the two-hybrid assay can then be determined by the culture of mammalian cells, for example human cells transfected with the xcex2TrCp protein or a fragment thereof, in the presence of a reporter gene specific for the protein involved in the pathological condition which it is desired to treat.\nThus, for the IxcexaB protein, it will be possible to use human cells originating from the cell lines Hela, 293, 293T, etc. and the reporter gene dependent on NFxcexaB sites (3Enh-xcexaB-ConA Luc), which controls the expression of luciferase.\nIn non-stimulated human cells, the human xcex2TrCP protein is transitorily expressed from a eukaryotic expression vector such as pCDNA3 (Invitrogen), or any other eukaryotic expression vector, which has inserted the DNA coding for the xcex2TrCP protein under the control of a strong promoter of the cytomegalovirus, CMV, or the like. An amount of the order of 3 xcexcg of this vector permitting the expression of the xcex2TrCP protein will be cotransfected by one of the common transfection techniques (calcium phosphate, lipofectamine (Life Technologies), electroporation (Ausubel and Sambrook, cf. below) etc.) with 1 xcexcg of a reporter vector dependent on NFxcexaB sites (3Enh-xcexaB-ConA luc) or independent of NFxcexaB sites (RSV Luc or ConA Luc) which control the expression of the luciferase reporter gene. Molecules capable of inhibiting the h-xcex2TrCP/IxcexaB interaction will inhibit the increase in the expression of lucifersse in this assay. These inhibitors will be added to the culture medium for at least 6 hours, 24, 36 or 48 hours after transfection. The specificity of these inhibitors may be checked by verifying that they have no effect on RSV Luc or ConA Luc. Another possible alternative will be to use the dual luciferase system from Promega, in which two different reporter vectors can be assayed at the same time.\nAccording to one experimental protocol similar to that described above, but with stimulated cells, it will be possible to verify that the inhibition induced by the expression of the h-xcex2TrCPxcex94F fragment on the TNF-dependent transcription activation has been nullified.\nThus, in this second assay, the human cells are cotransfected with 1 xcexcg of reporter vector, i.e. either 3Enh-xcexaB-ConA Luc, ConA Luc or RSV Luc, and with 3 xcexcg of pCDNA3 expressing the h-xcex2TrCPxcex94F peptide fragment, which is a mutant of xcex2TrCP with its F-box deleted. 24 to 48 h after transfection, the cells are treated for 6 h with TNF or okadaic acid (OKA), which are potent NFxcexaB activators (BAUERLE et al., Cell, 1996, 87, 13-20). The h-xcex2TrCPxcex94F mutant has a massive inhibitory effect on the expression of the luciferase reporter compared with a control plasmid transfected under the same conditions. This effect is due to the inhibition of IxcexaB degradation induced by the binding of the h-xcex2TrCPxcex94F mutant in place of the endogenous wild-type h-xcex2TrCP protein. Therefore an inhibitor of the h-xcex2TrCP/IxcexaB interaction will also inhibit the h-xcex2TrCPxcex94F/IxcexaB interaction and hence will reverse the inhibitory effect of the h-xcex2TrCPxcex94F fragment. The potential inhibitors are added to the medium under the same conditions as those indicated above. From the cells stimulated with TNF or OKA, those inhibitors are chosen which induce an increase in the expression of the reporter gene.\nAfter the selection of inhibitors in the previous two assays, a third assay can be carried out to verify that they are capable of inhibiting the activation of NFxcexaB induced by stimulation of the cells with TNF or OKA.\nThe cells transfected only with 1 ng of reporter vector (3Enh-KB-ConA Luc) and stimulated for 6 h with TNF or OKA are treated with the potential inhibitors. To be specific, these inhibitors must have an effect only on the IxcexaB-dependent reporter vectors and not on the other reporter vectors (ConA or RSV).\nIn the case of xcex2-catenin, it will be possible to use human cells originating from the above lines transformed with mutated xcex2-catenin or the peptide fragment of xcex2TrCP devoid of the F-box, in the presence of vector Top-TK-Luci, which contains a multimer of TCF-LEF sites responding to xcex2-catenin, patenin, or Fop-tk Luci, which contains an inactive mutated multimer and no longer responds to xcex2-catenin.\nFurthermore, as oncogenic mutated xcex2-catenin can easily be distinguished from wild-type xcex2-catenin by the fact that the former, in contrast to the latter, is incapable of binding to xcex2TrCP in the two-hybrid assay, xcex2-catenin mutations can be detected in human tumors by measuring the interaction with xcex2TrCP in the two-hybrid assay.\nThis assay is valuable because xcex2-catenin mutations are found in numerous cancers such as colon cancer, melanomas, hepatocarcinomas, etc. The only way of detecting these mutations hitherto was to sequence the xcex2-catenin by carrying out RT-PCR on the RNA of the tumors studied. For greater reliability, several double-stranded sequences have to be made in this assay of the prior art. Also, the existence of a mutation does not in itself indicate the oncogenic character of this mutation. It could be a case of polymorphism unconnected with tumorigenicity.\nThe advantage of the two-hybrid assay with the xcex2TrCP protein is that, in times equivalent to those required to obtain a sequence, it is possible to obtain a clear answer regarding the percentage of oncogenic mutated xcex2-catenin sequences detected from the tumoral RNA. Over a large number of colonies, the percentage of oncogenic forms of xcex2-catenin which are incapable of interacting with xcex2TrCP, compared with the wild-type forms which do interact with xcex2TrCP, can be determined precisely. The assay can be performed in a time equivalent to that required to obtain a few sequences, and at a reduced cost.\nThis assay comprises the following steps:\n1xe2x80x94Preparation of the total RNA from a biopsy of a tumor and of the surrounding healthy tissue, as control, using one of the various RNA preparation techniques or kits (AUSUBEL et al., Current Protocols in Molecular Biology).\n2xe2x80x94Amplification of the catenin sequences of the tumor and of the surrounding healthy tissue by carrying out RT-PCR on the RNA samples using a pair of oligonucleotides which permit amplification either of the N-terminal part only (1-130), which contains the most frequently encountered oncogenic mutations (RUBINFELD B. et al., Science, 275, 1790-1792, 1997; DE LA COSTE et al., Proc. Natl. Acad. Sci. USA, 95, 8847-8851, 1998), or of the whole of the xcex2-catenin coding sequence.\n3xe2x80x94Insertion of these amplified fragments, by ligation, into one of the two-hybrid vectors, for example pGAD1318, to give a frame fusion with the Gal4 transcription activation domain or the equivalent activation domain for transcription or binding to the DNA coded for by the two-hybrid vector used.\n4xe2x80x94ransformation of bacteria of various appropriate strains and plating of the whole of the transformant on LB-ampicillin medium.\n5xe2x80x94Harvesting of all the colonies and plasmid minipreparation (AUSUBEL, supra).\n6xe2x80x94L40 yeasts or any other appropriate strain of yeast will be cotransformed by the plasmid containing the xcex2-catenin sequences of the above minipreparation with a fusion hybrid containing xcex2TrCP, for example pLexA-xcex2TrCP, in which the xcex2TrCP is fused to the LexA DNA binding domain. A two-hybrid assay is performed on all the colonies obtained, for example by plating the cotransformed yeasts on DO-W-L medium and then transferring the colonies to selective medium for detection of the interactions, i.e. DOW-L-H medium, or in the presence of X-Gal for detection of the interactions by xcex2-galactosidase production (BARTEL P. and FIELDS S., Meth. Enzymol., 254, 241-263, 1995).\nThis assay requires the following reagents:\n1xe2x80x94Vector pGAD1318 predigested at the appropriate sites for inserting the amplified fragment obtained by RT-PCR.\n2xe2x80x94The appropriate oligonucleotides for amplifying the xcex2-catenin sequence and then inserting the amplified xcex2-catenin sequences. The oligonucleotide primers for amplification will be chosen according to the mode of insertion of the amplified fragment and the chosen sites.\n3xe2x80x94Plasmid pBTM116-xcex2TrCP expressing xcex2TrCP fused to the LexA DNA binding domain.\n4xe2x80x94As control, plasmids coding for fusion hybrids with control proteins, for example pLexRas and pGAD1318Raf.\nThe gap repair technique (SCHWARTZ H. et al., Mutation detection by a two-hybrid assay, Hum. Mol. Gen., 7, 1029-1032, 1998) may also be applied for this assay in order to insert the sequence of the amplified xcex2-catenin fragment into the two-hybrid vector and transform yeasts directly without proceeding via the step of prior transformation in bacteria.\nThe invention will now be described in detail with the aid of the following account of experiments.\nA large part of the techniques described in these Examples, which are well known to those skilled in the art, is explained in detail in the book by SAMBROOK et al. (supra) or in the book by AUSUBEL et al. (supra)."}
{"text": "Embodiments of the inventive subject matter generally relate to the field of wireless network devices and, more particularly, to symbol error detection for Bluetooth® basic data rate packets.\nBluetooth basic data rate (BDR) packets are modulated using Gaussian frequency shift keying (GFSK) modulation techniques. In addition to Bluetooth BDR packets, GFSK is also used to modulate a header of a Bluetooth enhanced data rate (EDR) packet. In a frequency shift keying (FSK) modulation scheme, a binary one is represented by a positive frequency deviation and a binary zero is represented by a negative frequency deviation. GFSK is a type of FSK modulation technique that utilizes a Gaussian filter to smooth positive/negative frequency deviations. The Gaussian pulse shaping improves spectral efficiency for Bluetooth signals."}
{"text": "1. Field of the Invention\nThe present invention relates to a mobile communication system and, particularly, to a mobile communication terminal and a communication method therefor which enable two-way communication with a base station.\n2. Description of the Prior Art\nIn general, mobile terminals carried by users are battery-powered. Therefore, it is desirable that their power consumption be minimized. There has been proposed a radio telephone terminal that is intended for such reduction in power consumption in Japanese Patent Application Laid-Open No. 7-303077. This radio telephone terminal is equipped with a detector for detecting the remaining amount of battery power. If the battery remaining power is greater than or equal to a predetermined amount, location registration is performed every time the telephone terminal moves to an adjacent communicable area (cell). If the battery remaining power is smaller than the predetermined amount, the user is merely informed of non-effectiveness of location registration. Further, the time interval of reception of control signals from a base station and that of location registration are varied in accordance with the battery remaining power, to reduce the battery consumption and in turn elongate the usable time.\nHowever, a transmission operation for location registration requires a large amount of power because of driving of a high-frequency amplifier. In the conventional radio communication terminal as described above, because it is necessary to perform, on a regular basis, transmission and reception operations for location registration with the base station, the battery is consumed at a high rate and the usable time cannot be made sufficiently long."}
{"text": "The present invention relates to an ultrasonic microscope, and more particularly to an ultrasonic microscope in which a specimen to be observed by the ultrasonic microscope can be observed by an optical microscope incorporated therein.\nIn recent years, it has become possible to generate and detect an ultrasonic wave having an ultrahigh frequency as high as 1 GHz and hence to realize an ultrasonic wavelength of about 1.5 .mu.m in the water. As a result, an ultrasonic imaging equipment having a high resolving power has been proposed (see, for example, U.S Pat. No. 4,386,530). FIG. 1 shows the schematic construction of an ultrasonic microscope. Referring to the figure, the transmission, focusing and reception of an ultrasonic acoustic wave is made by an acoustic lens 1 including a cylindrical block of an acoustic propagating medium such as fused quartz one end face of which is polished to provide an optically flat face. On the polished face is disposed a laminated structure which includes upper and lower electrodes 3 with a piezoelectric thin film 2 of zinc oxide or the like sandwiched therebetween. An electric pulse signal 5 from a pulse oscillator 4 is applied to the piezoelectric thin film 2 to generate an ultrasonic acoustic wave 6 therefrom. The other end face of the acoustic lens 1 is provided with a semispherical recess having a diameter of 0.1 to 1.0 mm, and a space between the semispherical recess portion and a specimen 7 is filled with a medium 8 (for example, water) for propagating the ultrasonic wave 6 to the specimen 7. With such a construction, the ultrasonic wave 6 generated from the piezoelectric thin film 2 is propagated through the cylindrical block 1 in the form of a plane wave and is refracted at the semispherical recess portion in accordance with a difference in sound velocity between the fused quartz and the medium 8 so that the plane wave is focused on the specimen 7. The ultrasonic wave reflected from the specimen 7 is collected and converted into a plane wave by the acoustic lens 1. The plane wave is propagated to the piezoelectric thin film 2 to be converted into a radio frequency (RF) signal 9. The RF signal 9 is received by a receiver 10 which includes a diode detector for conversion into a video signal 11. The video signal 11 is used as the input signal of a CRT display unit 12. The specimen 7 is subjected to a two-dimensional scanning in an X-Y plane on the basis of a signal from a power source 13 for driving a specimen table 14. The variation of the intensity of the reflected ultrasonic wave from the specimen 7 thus scanned is two-dimensionally displayed on the screen of the CRT display unit 12.\nThe acoustic wave used in an ultrasonic microscope has an ultra-high frequency and the attenuation of the acoustic wave in a medium such as water increases in proportion to the second power of the frequency if the transmitting distance is constant. Therefore, it is required to maintain a gap between the acoustic lens 1 and the specimen 7 as small as possible in order to make the transmitting distance short. For example, in the case of an ultrasonic wave having a frequency of 1 GHz, the gap between the acoustic lens 1 and the specimen 7 is selected to be about 50 .mu.m, thereby minimizing the attenuation of the ultrasonic wave. On the other hand, if the semispherical recess portion of the acoustic lens 1 fabricated with high precision is damaged by collision with the specimen, etc. during an observation operation as may take place in an optical lens, the damaged acoustic lens is rendered useless. The above-mentioned requirements of maintaining the gap between the acoustic lens 1 and the specimen 7 very small makes it very difficult to visually ascertain the presence of the gap. The requirements involve a very high possibility that the acoustic lens 1 may be damaged by collision with the specimen 7 during a manipulation of the lens."}
{"text": "In optical communication devices, an optical coupling lens is employed to couple optical signals between an optical emitter/receiver and an optical fiber. The optical coupling lens includes a light incident surface, a light emergent surface, and a reflecting surface. The reflecting surface is for reflecting optical signals between the light incident surface and the light emergent surface. The optical coupling lens includes at least one first converging lens formed on the light incident surface, and at least one second converging lens formed on the light emergent surface."}
{"text": "1. Field of the Invention\nThe present invention relates in general to an ocular lens such as a contact lens placed on an eyeball or an intraocular lens inserted within the eye. More particularly, the invention is concerned with a multifocal ocular lens having a plurality of vision correction regions with different optical or vision correction powers.\n2. Discussion of the Related Art\nAs an ocular lens used for vision correction of an eye having weakened accommodation faculty such as presbyopia, there has been proposed a bifocal or multifocal ocular lens having a plurality of vision correction regions which provide respective different optical powers. JP-A-63-95415 and JP-A-1-319729 disclose an alternating or translating vision type contact lens as one example of the multifocal ocular lens, wherein the vision correction regions with different optical powers are selectively and alternatively used as needed with a shift of the visual axis of the lens wearer. U.S. Pat. No. 4,580,882 and JP-A-2-217818 disclose a simultaneous vision type contact lens as another example of the bifocal or multifocal ocular lens, wherein the vision correction regions with different optical powers are simultaneously used, and a desired image observed through one of the vision correction regions is selected by an action of the wearer's brain.\nIn either of the alternating and simultaneous vision types described above, the vision correction regions of the bifocal lens consist of a near vision correction region through which near objects are observed, and a distant vision correction region through which distant objects are observed. On the other hand, the vision correction regions of the multifocal lens consist of such near and distant vision correction regions, and at least one intermediate region which is interposed between those near and distant vision correction regions and which provides an optical power different from the optical powers of the near and distant vision correction regions.\nThe bifocal ocular lens having the two vision correction regions, i.e., the near and distant vision correction regions, tends to suffer from a jump or overlapping of the image of an object located intermediate between near and far objects. This is because the bifocal ocular lens has only two focal points which respectively correspond to the near and distant vision correction regions. In this case, the obtained image is undesirably blurred. Further, the bifocal lens has surface discontinuity at a junction of the near and distant vision correction regions, in other words, the two vision correction regions are not smoothly connected to each other at the junction therebetween, deteriorating a wearing comfort as felt by the user.\nAs compared with the bifocal lens described above, the multifocal ocular lens having the intermediate region(s) between the near and distant vision correction regions provides a clearer image for the object located intermediate between the near and far objects. In the multifocal ocular lens, however, each of the plurality of vision correction regions inevitably has a small radial dimension, so that the optical power of the lens changes in steps in its radial direction, corresponding to the plurality of vision correction regions. In this case, the near and distant vision correction regions do not have sufficiently large surface areas, whereby the images to be obtained through these near and distant vision correction regions undesirably tend to be blurred. Moreover, the optical power in the multifocal lens changes in steps abruptly at each boundary between the adjacent two vision correction regions, in other words, the surfaces of the plurality of vision correction regions are not smoothly connected to one another with a continuous change of the optical power. Therefore, the multifocal lens also suffers from deteriorated wearing comfort. Further, the multifocal ocular lens tends to suffer from problems of so-called \"ghosting\" (ghost images) or double imaging, and mutual interference of the images obtained through the near and distant vision correction regions.\nThe assignee of the present invention proposed in JP-A-5-181096 a multifocal ocular lens having a near vision correction region, an intermediate region, and a distant vision correction region, which are located concentrically or coaxially with one another. In the proposed lens, the optical power continuously changes along a suitable curve from the value of one of the near and distant vision correction regions to the value of the other region. According to this arrangement, the optical power changes smoothly even at the junction of the near vision correction region and the intermediate region, and at the junction of the intermediate region and the distant vision correction region, and the lens has a smooth surface configuration which assures the user of a comfortable wearing, without having any surface discontinuity at the junctions. Further, the proposed multifocal ocular lens is free from the problem of ghosting experienced in the conventional multifocal lens wherein the vision correction regions are not smoothly connected to one another.\nAs a result of an extensive study by the inventors of the present invention, it was found that even the proposed multifocal ocular lens does not sufficiently meet the users' requirements. Described in detail, the proposed lens is still unsatisfactory in providing sufficiently clear viewing of both of the near and distant objects when the lens is used in specific conditions or environments by a painter and a surveying engineer, for instance."}
{"text": "1. Field of the Invention\nThe field of invention relates to cleat guard apparatus, and more particularly pertains to a new and improved shoe cleat guard wherein the same is arranged to afford protection to underlying surfaces in mounting the guard structure to an associated athletic shoe having cleats and spikes.\n2. Description of the Prior Art\nCleat guard structure of various types is available in the prior art and exemplified by the U.S. Pat. Nos. 4,872,273; 5,070,631; and 4,484,398 as examples, wherein the instant invention attempts to overcome deficiencies of the prior art by providing for a cleat guard structure arranged to be pierced by cleats of an associated shoe to insure securement and positioning of the cleats in a non-slip relationship relative to the guard structure and in this respect, the present invention substantially fulfills this need."}
{"text": "1. Field of the Invention\nThe present invention relates to a head slider supporting device, a disk device and a suspension, and, in particular, to a magnetic head slider supporting device, an optical head slider supporting device, a magnetic disk device, an optical disk device and a suspension.\nAs a result of an increase of the frequency of a signal which an information processing apparatus processes, it is requested that the write current frequency of a magnetic disk device be increased from 70 MHz to 200 through 300 MHz, for example. In order to increase the write current frequency, it is necessary to reduce the inductance and the electrostatic capacity of a signal transmission path from a magnetic head slider to a head IC. For this purpose, it is effective to provide the head IC at a position near to the magnetic disk slider. Further, thinning of the magnetic disk device is also requested. It is necessary that the head IC be loaded in the magnetic disk device in a condition in which the thinning of the magnetic disk device is not disturbed, and, also, that the head IC does not come into contact with a magnetic disk and so forth even when a shock is applied to the magnetic disk device.\n2. Description of the Related Art\nMagnetic disk devices, in each of which a head IC for amplifying a signal read through a magnetic head slider is installed on an arm, are disclosed in Japanese Laid-Open Patent Application Nos.62-217476, 3-108120, 3-187295, 3-192513 and so forth.\nHowever, in each of the above-mentioned magnetic disk devices, because the distance between the magnetic head slider and the head IC is long, it is difficult to reduce the inductance and the electrostatic capacity of the signal transmission path from the magnetic head slider to the head IC. Further, the head IC is packaged in a synthetic resin, and, thereby, is thick. As a result, in order to prevent the head IC from coming into contact with a magnetic disk and so forth when a shock is applied to the magnetic disk device, it is necessary to elongate the distance between adjacent magnetic disks. As a result, the magnetic disk device is thick. Further, because the head IC is packaged in the synthetic resin, the head IC is heavy. As a result, the magnetic head slider is heavy. Thereby, the flying stability of the magnetic head slider above the magnetic disk is degraded, and also, there is a possibility that, when the magnetic head slider comes into contact with a magnetic disk due to a strong shock applied to the magnetic disk device, the shock applied to the magnetic disk is so strong that the magnetic disk is damaged.\nAs shown in FIG. 1, in a head slider supporting device 1, on the top surface 2a of a suspension 2 (hereinafter, the position of the suspension shown in FIG. 1 is a reference position thereof, and xe2x80x98the top surfacexe2x80x99 of the suspension means the top surface in this position of the suspension), wiring patterns 3 are formed from the extending end to the fixed end of the suspension 2, and the magnetic head slider 4 is loaded on the top surface 2a of the suspension 2 at the extending end of the suspension 2.\nHere, provision of the head IC 5 will be considered. Due to the arrangement of the wiring patterns 3, a surface at which the head IC is loaded is limited to the top surface 2a of the suspension 2. When considering an increase of the write current frequency, it is preferable that the head IC 5 be provided at a position near to the magnetic head slider 4. Therefore, it is assumed that the head IC 5 is loaded on the top surface 2a of the suspension 2 near the magnetic head slider 4.\nIn order to prevent the head IC 5 from coming into contact with a magnetic disk 6 even when a strong shock is applied to the magnetic disk device, it is necessary that a gap 7 of the distance xe2x80x98axe2x80x99 equal to or longer than 0.15 mm to be provided between the head IC 5 and the magnetic disk 6.\nRecently, in order tot thin the magnetic disk device, the magnetic head slider 5 of a small size (a so-called pico-slider, the height xe2x80x98bxe2x80x99 of which is 0.3 mm) has been used. As a result, the distance xe2x80x98cxe2x80x99 between the suspension 2 and the magnetic disk 6 is small.\nWhen considering a bare head IC 5, the bare head IC 5 is cut out from a wafer. Accordingly, the thickness of the bare IC 5 is determined by the thickness of the wafer. At the present time, it is difficult to thin the wafer to less than 0.3 mm. Accordingly, the height (thickness) xe2x80x98dxe2x80x99 of the head IC 5 is approximately 0.3 mm minimum.\nTherefore, when the bare head IC 5 is simply loaded on the top surface 2a of the suspension 2, it is difficult to obtain the gap equal to or longer than 0.15 mm between the head IC 5 and the magnetic disk 6. Thus, a special device is needed when the bare head IC 5 is loaded on the top surface 2a of the suspension 2.\nAn object of the present invention is to provided a head slider supporting device, a disk device and a suspension by which the above-described problem is solved.\nA head slider supporting device, according to the present invention, comprises:\na head slider having a head loaded thereon;\na head IC which controls the head;\na substantially flat supporting member having a portion at which the head slider is loaded and a head IC mounting portion at which the head IC is mounted; and\nwires electrically connecting the head slider and the IC,\nwherein the head IC is mounted at the head IC mounting portion in a condition in which the head IC projects from one surface of the supporting member so that an amount of projecting of the head IC from the one surface of the supporting member is smaller than the thickness of the head IC.\nThus, the length by which the head IC protrudes from the surface of the supporting member is shorter than the thickness of the head IC. As a result, in a case where a so-called pico-slider is used as the head slider, a gap can be provided between the head IC and a disk such that, even when a strong shock is applied to the disk device, the head IC is prevented from coming into contact with the disk.\nA head slider supporting device, according to another aspect of the present invention, comprises;\na suspension, a head IC chip mounting portion being provided on a first surface of the suspension;\na head slider loaded on the first surface of the suspension at an extending end thereof, the head slider integrally including a head;\na head IC chip mounted at the head IC chip mounting portion of the suspension; and\nwiring patterns extending along the suspension from a portion of the suspension, at which portion the head slider is loaded, via the head IC chip mounting portion,\nwherein the head IC chip is mounted at the head IC chip mounting portion in a condition in which the head IC chip is lowered so that an amount of projecting of the head IC chip from the first surface of the suspension is smaller than the thickness of the head IC chip.\nThus, the length by which the head IC chip protrudes from the first surface of the suspension is shorter than the thickness of the head IC chip. As a result, in a case where a so-called pico-slider is used as the head slider, a gap can be provided between the head IC chip and a disk such that, even when a strong shock is applied to the disk device, the head IC chip is prevented from coming into contact with the disk. Thus, the disk device can be provided in which the pico-slider is used as the head slider, and also, the head IC chip is mounted on the first surface of the suspension on which the head slider is loaded. Each of the wiring patterns, electrically connecting the head slider with the head IC chip, can be as short as several millimeters, because the head IC chip is mounted on the first surface on which the head slider is loaded. As a result, the inductances of the wiring patterns are small. Further, the electrostatic capacities between adjacent wiring patterns are small. Therefore, in a case where the disk device is the magnetic disk device, a signal of, for example, 200 MHz, higher than 70 MHz, as in the case of the related art, can be written in and read from the magnetic disk in the magnetic disk device.\nThe head IC chip mounting portion may extend in the suspension along the longitudinal direction of the suspension, both sides of the head IC chip mounting portion being cut and a middle portion of the head IC chip mounting portion projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on the middle portion of the head IC chip mounting portion.\nThereby, the head IC chip mounting portion can be formed without drawing the suspension. Accordingly, the head IC chip mounting portion can be formed without an excessive stress being applied to the suspension. Further, because both sides of the head IC chip mounting portion are cut along the longitudinal direction of the suspension, the head IC chip mounting portion can be formed with the characteristics of the suspension being affected thereby as little as possible.\nThe head IC chip may comprise a first half portion and a second half portion, the first half portion having an extending portion, which extends laterally beyond the second half portion of the head IC chip on the side of the first half portion; and\nthe head IC chip mounting portion includes an opening having a size such that the second half portion of the head IC chip passes through the opening and the extending portion of the first half portion of the head IC chip is supported by a peripheral portion of the opening.\nThereby, the head IC chip mounting portion can be formed with the characteristics of the suspension being affected thereby as little as possible. Further, merely by causing the second half portion of the head IC chip to pass through the opening and causing the extending portion of the first half portion of the head IC chip to be supported by the peripheral portion of the opening, the height by which the head IC chip protrudes from the first surface of the suspension can be precisely determined. Further, by determining the size of the opening such that the second half portion fits into the opening, providing terminals on the bottom surface of the head IC chip and providing terminals of the wiring patterns along the periphery of the opening, it is possible that the terminals of the head IC chip precisely face the terminals of the wiring patterns merely by causing the second half portion of the head IC chip to pass through the opening. Thereby, it is possible to electrically connect the terminals of the head IC chip with the terminals of the wiring patterns with high reliability.\nThe head IC chip mounting portion may comprise an opening having a size such that a certain portion of the head IC chip passes through the opening.\nThereby, the head IC chip mounting portion can be formed with the characteristics of the suspension being affected thereby as little as possible. In this case, terminals may be provided on side walls of the head IC chip, terminals of the wiring patterns may be provided on a peripheral portion of the opening, and the terminals on the side walls of the head IC chip may be electrically connected with the terminal of the wiring patterns. Further, it is possible to mount the head IC chip at the head IC chip mounting portion in a condition in which a portion of the head IC chip is lowered from the first surface of the suspension.\nThe suspension may have a rigid portion between the extending end and the other end, the rigid portion having a rib on at least one side thereof so that the rigid portion is prevented from bending; and\nthe head IC chip mounting portion may be formed in the rigid portion.\nThereby, the head IC chip mounting portion can be formed with the characteristics of the suspension being affected thereby as little as possible.\nA disk device, according to the present invention, comprises:\nan actuator;\na disk which is rotated;\nan arm which is driven by the actuator; and\na head slider supporting device, which is rotated integrally with the arm, the head slider supporting device comprising:\na head slider having a head loaded thereon;\na head IC which controls the head;\na substantially flat supporting member having a portion at which the head slider is mounted and a head IC mounting portion at which the head IC is mounted; and\nwires electrically connecting the head slider and the head IC,\nwherein the head IC is mounted at the head IC mounting portion in a condition in which the head IC projects from one surface of the supporting member so that an amount of projecting of the head IC from the one surface of the supporting member is smaller than the thickness of the head IC.\nA disk device, according to another aspect of the present invention, comprising:\nan actuator;\na disk which is rotated;\nan arm which is driven by the actuator; and\na head slider supporting device, which is rotated integrally with the arm, the head slider supporting device comprising:\na suspension, a head IC chip mounting portion being provided on one surface of the suspension;\na head slider loaded on the one surface of the suspension at an extending end thereof, the head slider integrally including a head;\na head IC chip mounted at the head IC chip mounting portion of the suspension; and\nwiring patterns extending along the suspension from a portion of the suspension, at which portion the head slider is loaded, via the head IC chip mounting portion,\nwherein the head IC chip is mounted at the head IC chip mounting portion in a condition in which the head IC chip is lowered so that an amount of projecting of the head IC chip from the one surface of the suspension is smaller than the thickness of the head IC chip.\nThus, as described above, the length by which the head IC chip protrudes from the one surface of the suspension is shorter than the thickness of the head IC chip. As a result, in a case where a so-called pico-slider is used as the head slider, a gap can be provided between the head IC chip and a disk such that, even when a strong shock is applied to the disk device, the head IC chip is prevented from coming into contact with the disk. Thus, the disk device can bed provided in which the pico-slider is used as the head slider, and, also, the head IC chip is mounted on the one surface of the suspension. Each of the wiring patterns electrically connecting the head slider with the head IC chip, can be short, because the head IC chip is mounted on the one surface on which the head slider is loaded. As a result, the inductances of the first wiring patterns are small. Further, the capacities between adjacent first wiring patterns are small. Therefore, in a case where the disk device is the magnetic disk device, a signal of, for example, 200 MHz, higher than 70 MHz, as in the case of the related art, can be written in and read from the magnetic disk in the magnetic disk device.\nA suspension, according to the present invention, has an extending end and the other end, a head IC chip mounting portion being provided between the extending end and the other end thereof, wherein:\na head slider is loaded on the suspension at the extending end thereof, the head slider integrally including a head;\na head IC chip is mounted at the head IC chip mounting portion of the suspension;\nfirst wiring patterns extend along the suspension between a portion of the suspension, at which portion the head slider is loaded, and the head IC chip mounting portion; and\nsecond wiring patterns extend along the suspension from the head IC chip mounting portion,\nwherein the head IC chip is mounted at the head IC chip mounting portion in a condition in which the head IC chip projects from one surface of the suspension so that an amount of projecting of the head IC chip from the one surface of the suspension is smaller than the thickness of the head IC chip.\nThus, as described above, the length by which the head IC chip protrudes from the one surface of the suspension is shorter than the thickness of the head IC chip. As a result, in a case where a so-called pico-slider is used as the head slider, a gap can be provided between the head slider and a disk such that, even when a strong shock is applied to the disk device, the head IC chip is prevented from coming into contact with the disk. Thus, the disk device can be provided in which the pico-slider is used as the head slider. The head slider can be loaded on the one surface of the suspension on which the head IC chip is mounted. Each of the first wiring patterns, electrically connecting the head slider with the head IC chip, can be short, because the head IC chip is mounted on the one surface on which the head slider is loaded. As a result, the inductances of the first wiring patterns are small. Further, the electrostatic capacities between adjacent first wiring patterns are small. Therefore, in a case where the disk device is the magnectic disk device, a signal of, for example, 200 MHz, higher than 70 MHz, as in the case of the related art, can be written in and read from the magnetic disk in the magnetic disk device.\nThe head IC chip mounting portion may extend in the suspension along the longitudinal direction of the suspension, both sides of the head IC chip mounting portion being cut and a middle portion of the head IC chip mounting portion projecting on the side of the other surface of the suspension, the head IC chip being mounted on the middle portion of the head IC chip mounting portion.\nThereby, the head IC chip mounting portion can be formed without drawing the suspension. Accordingly, the head IC chip mounting portion is formed without an excessive stress being applied to the suspension. Further, because both sides of the head IC chip mounting portion are cut along the longitudinal direction of the suspension, the head IC chip mounting portion can be formed with the characteristics of the suspension being affected thereby as little as possible.\nThe head IC chip mounting portion may comprise an opening having a size such that a certain portion of the head IC chip passes through the opening.\nThereby, the head IC chip mounting portion can be formed with he characteristics of the suspension being affected thereby as little as possible.\nThe suspension may have a rigid portion having a rib portion on at least one side of the rigid portion so that the rigid portion is prevented from bending; and\nthe head IC chip mounting portion may be formed in the rigid portion.\nThereby, the head IC chip mounting portion can be formed with the characteristics of the suspension being affected thereby as little as possible.\nBoth sides of the head IC chip mounting portion may be cut, and, also, a center of the head IC chip mounting portion may be cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that the head IC chip mounting portion is divided into two parts, middle portions projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on, so as to extend over, the middle portions of the two parts of the head IC chip mounting portion on surfaces thereof on the side of the first surface of the suspension; and\nthe wiring patterns reach the middle portions via slope portions which are formed at both ends of the head IC chip mounting portion when the middle portions project on the side of the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portions.\nIn this arrangement, the slope portions are merely bent but are not lengthened when a press machine is used for causing the middle portion of the head IC chip mounting portion to project on the side of the second surface. Thereby, cutoff of the wiring patterns formed on the slope portions can be effectively avoided.\nBoth sides of the head IC chip mounting portion may be cut, and, also, one end of the head IC chip mounting portion may be cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that an opening is formed but a thin connection portion is left which bridges the opening, a middle portion of the head IC chip mounting portion projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on the middle portion of the head IC chip mounting portion on a surface thereof on the side of the first surface of the suspension; and\nthe wiring patterns reach the middle portion via a slope portion which is formed at the other end of the head IC chip mounting portion when the middle portion projects from the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portion.\nIn this arrangement, the thin connection portion effectively lengthens while the slope portion is only bent but is not lengthened. Thereby, cutoff of the wiring patterns formed on the slope portion can be effectively avoided.\nA head slider supporting device, according to another aspect of the present invention, comprises:\na suspension, having an extending end and the other end, having a fixing portion on the side of the other end, and having a head IC chip mounting portion at the fixing portion;\na head slider loaded on a first surface of the suspension at the extending end thereof, the head slider integrally including a head;\na head IC chip mounted at the head IC chip mounting portion of the suspension;\nwiring patterns extending along the suspension from a portion of the suspension, at which portion the head slider is loaded, to the head IC chip mounting portion, and further extending from the head IC chip mounting portion; and\na plate-shaped junction member, on which the fixing portion of the suspension is mounted,\nwherein:\nboth sides of the head IC chip mounting portion are cut, and, also, a center of the head IC chip mounting portion is cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that the head IC chip mounting portion is divided into two parts, and middle portions of the two parts of the head IC chip mounting portion projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on, so as to extend over, the middle portions of the two parts of the head IC chip mounting portion on surfaces thereof on the side of the first surface of the suspension;\nthe wiring patterns reach the middle portions via slope portions which are formed at both the ends of the head IC chip mounting portion when the middle portions project from the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portions; and\nthe junction member has an opening in which the head IC chip mounting portion at which the head IC chip is mounted is inserted.\nIn this arrangement, the slope portions are merely bent but are not lengthened when a press machine is used for causing the middle portion of the head IC chip mounting portion to project on the side of the second surface. Thereby, cutoff of the wiring patterns formed on the slope portions can be effectively avoided.\nFurther, as a result of the junction member having the opening, the head IC chip mounting portion is provided at the fixing portion located on the side of the other end of the suspension. That is, a place where the head IC chip is located is near to the rotation axis of the head slider supporting device. In comparison to a case where the head IC chip is located near to the extending end of the suspension, the inertia moment when the head slider supporting mechanism operates is small, and, thereby, it is possible to achieve a high-accuracy, high-speed seeking operation.\nA head slider supporting device, according to another aspect of the present invention, comprises:\na suspension, having an extending end and the other end, having a fixing portion at the other end, and having a head IC chip mounting portion at the fixing portion;\na head slider loaded on a first surface of the suspension at the extending end thereof, the head slider integrally including a head;\na head IC chip mounted at the head IC chip mounting portion of the suspension;\nwiring patterns extending along the suspension from a portion of the suspension, at which portion the head slider is loaded, to the head IC chip mounting portion, and further extending from the head IC chip mounting portion; and\na plate-shaped junction member, on which the fixing portion of the suspension is mounted,\nwherein:\nboth sides of the head IC chip mounting portion are cut, and, also, one end of the head IC chip mounting portion is cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that an opening is formed but a thin connection portion is left which bridges the opening, a middle portion of the head IC chip mounting portion projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on the middle portion of the head IC chip mounting portion on a surface thereof on the side of the first surface of the suspension;\nthe wiring patterns reach the middle portion via a slope portion which is formed at the other end of the head IC chip mounting portion when the middle portion projects from the second surface, and pad terminal which terminate the wiring patterns, respectively, are provided on the middle portion; and\nthe junction member has an opening, in which the head IC chip mounting portion at which the head IC chip is mounted is inserted.\nIn this arrangement, the thin connection portion effectively lengthens while the slope portion is only bent but is not lengthened. Thereby, cutoff of the wiring patterns formed on the slope portion can be effectively avoided.\nFurther as a result of the junction member having the opening, the head IC chip mounting portion is provided at the fixing portion located on the side of the other end of the suspension. That is, a place where the head IC chip is located is near to the rotation axis of the head slider supporting device. In comparison to a case where the head IC chip is located near to the extending end of the suspension, the inertia moment when the head slider supporting mechanism operates is small, and, thereby, it is possible to achieve a high-accuracy, high-speed seeking operation.\nA head slider supporting device, according to another aspect of the present invention, comprises:\na suspension, having a first surface, and having a head IC chip mounting portion on the second surface;\na head slider loaded on the first surface of the suspension at an extending end thereof, the head slider integrally including a head;\na head IC chip mounted at the head IC chip mounting portion of the suspension; and\nwiring patterns extending on the first surface of the suspension from a portion of the suspension, at which portion the head slider is loaded, via the head IC chip mounting portion;\nwherein:\nthe suspension includes a metal plate body and a base layer formed on the metal plate body, the wiring patterns being formed on the base layer;\nthe head IC chip mounting portion includes an opening formed in the metal plate body and having a size corresponding to the head IC chip, the base layer covering the opening on the side of the first surface, and pad terminals which terminate the wiring patterns, respectively, are exposed on a surface of the base layer on the side of the second surface of the suspension; and\nthe head IC chip is mounted at the head IC chip mounting portion so that the head IC chip is fitted into the opening of the metal plate body from the side of the second surface.\nIn this arrangement, as a result of the head IC chip being fitted into the opening from the side of the second surface. Further, the opening is formed in the suspension but the suspension is not bent. As a result, cutoff of the wiring patterns can be effectively avoided.\nA head slider supporting device, according to another aspect of the present invention, comprises:\na suspension, having an extending end and the other end, having a first surface and a second surface opposite to the first surface, having a fixing portion on the side of the other end, and having a head IC chip mounting portion at the fixing portion;\na head slider loaded on the first surface of the suspension at the extending end thereof, the head slider integrally including a head;\na head IC chip mounted at the head IC chip mounting portion of the suspension;\nwiring patterns extending on the first surface of the suspension from a portion of the suspension, at which portion the head slider is loaded, to the head IC chip mounting portion, and further extending from the head IC chip mounting portion; and\na plate-shaped junction member on which the fixing portion of the suspension is mounted,\nwherein:\nthe suspension includes a metal plate body and a base layer formed on the metal plate body, the wiring patterns being formed on he base layer;\nthe head IC chip mounting portion includes an opening formed in the metal plate body and having as size corresponding to the head IC chip, the base layer covering the opening on the side of the first surface, and pad terminals which terminate the wiring patterns, respectively, are exposed on a surface of the base layer on the side of the second surface of the suspension;\nthe head IC chip is mounted at the head IC chip mounting portion so that the head IC chip is fitted into the opening of the metal plate body from the side of the second surface; and\nthe junction member has an opening, in which the head IC chip mounted at the head IC chip mounting portion is inserted.\nIn this arrangement, as a result of the head IC chip being fitted into the opening from the side of the second surface, it is possible to shorten the length by which the head IC chip projects from the second surface. Further, the opening is formed in the suspension but the suspension is not bent. As a result, cutoff of the wiring patterns can be effectively avoided.\nFurther, as a result of the junction member having the opening, the head IC chip mounting portion is provided at the fixing portion located on the side of the other end of the suspension. That is, a place where the head IC chip is located is near to the rotation axis of the head slider supporting device. In comparison to a case where the head IC chip is located near to the extending end of the suspension, the inertia moment when the head slider supporting mechanism operates is small, and, thereby, it is possible to achieve a high-accuracy, high-speed seeking operation.\nA disk device, according to another aspect of the present invention, comprises:\nan actuator;\na disk which is rotated;\nan arm which is driven by the actuator; and\nany one of the above-described head slider supporting devices, which is rotated integrally with the arm.\nIn this arrangement, because any one of the above-described head slider supporting devices is used, the head IC chip is prevented from hitting the magnetic disk even when a strong shock is applied to the magnetic disk device. Further, it is possible to achieve the magnetic disk device which can write and read a signal of, for example, up to 200 MHz, higher than 70 MHz, as in the case of the related art.\nA suspension, according to another aspect of the present invention, has an extending end and the other end, has a fixing portion on the side of the other end, which portion is mounted on a junction member, and has a head IC chip mounting portion at the fixing portion,\nwherein:\na head slider is loaded on a first surface of the suspension at the extending end thereof, the head slider integrally including a head;\na head IC chip is mounted at the head IC chip mounting portion of the suspension; and\nwiring patterns extend along the suspension from a portion of the suspension, at which portion the head slider is loaded, to the head IC chip mounting portion, and further extend from the head IC chip mounting portion,\nwherein:\nboth sides of the head IC chip mounting portion are cut, and, also, a center of the head IC chip mounting portion is cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that the head IC chip mounting portion is divided into two parts, and middle portions of the two parts of the head IC chip mounting portion projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on surfaces of the middle portions of the two parts of the head IC chip mounting portion on the side of the first surface of the suspension; and\nthe wiring patterns reach the middle portions via slope portions which are formed at both ends of the head IC chip mounting portion when the middle portions project from the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portions.\nA suspension, according to another aspect of the present invention, has an extending end and the other end, has fixing portion on the side of the other end, which portion is mounted on a junction member, and has a head IC chip mounting portion at the fixing portion,\nwherein:\na head slider is loaded on a first surface of the suspension at the extending end thereof, the head slider integrally including a head;\na head IC chip is mounted at the head IC chip mounting portion of the suspension; and\nwiring patterns extend along the suspension from a portion of the suspension, at which portion the head slider is loaded, to the head IC chip mounting portion, and further extend from the head IC chip mounting portion,\nwherein:\nboth sides of the head IC chip mounting portion are cut, and, also, one end of the head IC chip mounting portion is cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that an opening is formed but a thin connection portion is left which bridges the opening, a middle portion of the head IC chip mounting portion projecting on the side of a second surface of the suspension, which second surface is opposite to the first surface, the head IC chip being mounted on the middle portion of the head IC mounting portion on a surface thereof on the side of the first surface of the suspension; and\nthe wiring patterns reach the middle portion via a slope portion which is formed at the other end of the head IC chip mounting portion when the middle portion projects from the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portion.\nA suspension, according to another aspect of the present invention, includes a metal plate body, a base layer formed on the metal plate body, and wiring patterns formed on the base layer, the suspension having an extending end and the other end, having a fixing portion on the side of the other end, which portion is mounted on a junction member, and having a head IC chip mounting portion at the fixing portion,\nwherein:\na head slider is loaded on a first surface of the suspension at an extending end thereof, the head slider integrally including a head;\na head IC chip is mounted at the head IC chip mounting portion of the suspension;\nwiring patterns extend on the first surface of the suspension from a portion of the suspension, at which portion the head slider is loaded, and further extend from the head IC chip mounting portion; and\nthe head IC chip mounting portion includes an opening formed in the metal plate body and having a size corresponding to the head IC chip, the base layer covering the opening on the side of the first surface, and pad terminals which terminate the wiring patterns, respectively, are exposed on a surface of the base layer on the side of a second surface of the suspension, which second surface is opposite to the first surface.\nIn each of these arrangements, because the head IC chip mounting portion is formed at the fixing portion, it is possible to achieve the suspension in which the head IC mounting portion is formed with the characteristics of the suspension being not affected thereby.\nA suspension, according to another aspect of the present invention, has a first surface and a second surface opposite to the firs surface, and has a rigid portion which has a rib portion on at least one side thereof so that the rigid portion is prevented from bending, a head IC chip mounting portion being formed in the rigid portion,\nwherein:\nthe head IC chip mounting portion has a head IC chip mounted thereon;\nboth sides of the head IC chip mounting portion are cut, and, also, a center of the head IC chip mounting portion is cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that the head IC chip mounting portion is divided into two parts, and middle portions of the two parts of the head IC chip mounting portion project on the side of the second surface of the suspension, the head IC chip being mounted on surfaces of the middle portions of the two parts of the head IC chip mounting portion on the side of the first surface of the suspension; and\nthe wiring patterns reach the middle portions via slope portions which are formed at both ends of the head IC chip mounting portion when the middle portions project on the side of the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portions.\nA suspension, according to another aspect of the present invention, has a first surface and a second surface opposite tot he first surface, and has a rigid portion which has a rib portion on at least one side thereof so that the rigid portion is prevented from bending, a head IC chip mounting portion being formed in the rigid portion,\nwherein:\nthe head IC chip mounting portion has a head IC chip mounted thereon;\nboth sides of the head IC chip mounting portion are cut, and, also, one end of the head IC chip mounting portion is cut in a direction perpendicular to a direction in which both sides of the head IC chip mounting portion are cut so that an opening is formed but a thin connection portion is left which bridges the opening, a middle portion of the head IC chip mounting portion projecting on the side of the second surface of the suspension, the head IC chip being mounted on the middle portion of the head IC mounting portion on a surface thereof on the side of the first surface of the suspension; and\nthe wiring patterns reaches the middle portion via a slope portion which is formed at the other end of the head IC chip mounting portion when the middle portion projects on the side of the second surface, and pad terminals which terminate the wiring patterns, respectively, are provided on the middle portions.\nA suspension, according to another aspect of the present invention, has a rigid portion which has a rib portion on at least one side thereof so that the rigid portion is prevented from bending, a head IC chip mounting portion being formed in the rigid portion,\nwherein:\nthe head IC chip mounting portion has a head IC chip mounted thereon;\nthe head IC chip mounting portion includes an opening formed in a metal plate body and having a size corresponding to the head IC chip, a base layer covering the opening on the side of the first surface of the suspension, pad terminals which terminate the wiring patterns, respectively, being exposed on a surface of the base layer on the side of a second surface of the suspension, which second surface is opposite to the first surface.\nIn each of these arrangements, because the head IC chip mounting portion is formed in the rigid portion which is prevented from bending, it is possible to achieve the suspension in which the head IC mounting portion is formed with the characteristics of the suspension being not affected thereby.\nOther objects and further features of the present invention will become more apparent from the following detailed description when read in conjunction with the accompanying drawings."}
{"text": "With rapid development of communication technology, particularly radio technology, a train running control system based on communication technology is researched increasingly deeply, and a Communication Based Train Control System (CBTC system) emerges accordingly. The CBTC system is a control system for realizing safety protection of the train based on the bidirectional train-to-ground wireless communication technology and the moving block control principle. The CBTC system is used for replacing the track circuit as a medium to realize running control of the train. The CBTC system is increasingly widely applied to train running safety protection of urban rail transit.\nOnly if the CBTC system recognizes the marshalling length of a train entering into a CBTC control area and determines that no vehicle (e.g., a track engineering maintenance vehicle) exists in a safety range in front of/in the rear of the train, the CBTC system provides movement authorization for the train. That is, the CBTC system transmits a safe running direction, a safe running distance and other information calculated accurately to the train. In the existing CBTC system, a transfer track is generally provided at the boundary (such as a connecting line between a departure line and other lines) of the CBTC control area, and a trackside CBTC device is provided beside the transfer track. The trackside CBTC device receives information reported by an on-board CBTC device of the train entering into the CBTC control area, and parses the received information to obtain the marshalling length of the train. The trackside CBTC device transmits movement authorization to the on-board CBTC device only in a case of determining the marshalling length of the train and determining that no vehicle exists in the safety range in front of/in the rear of the train.\nSince different trains usually have multiple different marshalling lengths such as six compartments or eight compartments, hybrid running of trains with different marshalling lengths is to be realized in a running line of the train. The existing CBTC system has to quickly determine whether a vehicle exists in a safety range in front of/in the rear of the train, in order to timely provide the movement authorization for the train. Therefore, only one type of train with the fixed marshalling length can run in the same running line generally, which limits the hybrid running of trains with different marshalling lengths."}
{"text": "There is a well-recognized need to clean-up contaminants that exist in ground water, i.e., aquifers and surrounding soil formations. Such aquifers and surrounding soil formations may be contaminated with various constituents including organic compounds such as, volatile hydrocarbons, including chlorinated hydrocarbons such as trichloroethene (TCE), and tetrachloroethene (PCE). Other contaminates that can be present include vinyl chloride, 1,1,1 trichloroethane (TCA), and very soluble gasoline additives such as methyltertiarybutylether (MTBE). Other contaminants may also be encountered.\nOzone sparging is now widely recognized as being one of the more effective oxidation techniques for destroying contaminants that exist in groundwater.\nOther types of contaminants are more recalcitrant. For instance, pharmaceuticals are particularly resistant to decomposition from known techniques including ozone sparging. Pharmaceuticals enter the groundwater from various sources. One source is pharmaceutical laboratories and manufacturing plants located in area with septic systems for waste disposal. Other sources are hospitals and nursing homes.\nPharmaceutical residuals are increasingly found in sewage discharges. Stronger selective oxidation techniques are necessary to the discharge of antibiotic-resistant bacteria and pharmaceutical residuals into groundwater and surface waters. Zwiener and Frimmel “Oxidative Treatment Of Pharmaceuticals In Water” Water Research 34(6) 1881-1885 (2000); Andreozzi et al. “Paracetanol Oxidation From Aqueous Solutions By Means Of Ozonation and H2O2 UV System” Water Research 37 993-1004 (2003), and Huber et al. “Oxidation Of Pharmaceuticals During Ozonation And Advanced Oxidation Processes” Environmental Science and Technology (2003) have proposed that oxidation systems need to be improved to address the variety of compounds involved. Korhonen et al. Oxidation of Selected Pharmaceuticals in Drinking Water Treatment, Presented at the ninth International Conference on Advanced Oxidation Technologies for Water and Air Remediation Canada (2003) felt that ozone or a combination of ozone and peroxide may offer effective treatment. The identified pharmaceutical residuals include the lipid regulator bezafibrate, antiepileplic carbamazepine, analgesic/inflammatory diclofenac and ibuprofen, and the antibiotic sulfamethoxazole. Even though Korhonen et al. (2003) obtained 90% removal of bezafibrate with ozone alone, H2O2 (peroxide) additional was necessary to obtain over 90% removal of carbamazepine, ibuprofen, and bezafibrate in clean water samples. However, with sewage, the presence of natural organic material (NOM) inhibits effective reaction.\nAnother need improved oxidation systems comes from treatment of alkanes and alkenes, common to petroleum products and spills. The bulk of petroleum products are aliphatic long-chain compounds, which are often 75% of the product. In heavier refined products, the carbon chain notation for molecular size, C5 to C30 denotes the dominant molecular fractions from 5-carbon to 30-carbon atoms strung together in a single chain. The higher fractions, particularly when branched, are resistant to bacterial action. Fogel (2001) has found that well-aerated samples of petroleum from a diesel source, even when supplied optimal nutrients, will leave about 25% undigested."}
{"text": "The present invention relates to an ink-jet print method and apparatus for printing lines by discharging an ink from an ink-jet head onto a recording member, a color filter, a display device, an apparatus having the display device, an ink-jet head unit adjusting device, and an ink-jet head unit.\nWith recent advances in personal computers, especially portable personal computers, the demand tends to arise for liquid crystal displays, especially color liquid crystal displays. However, in order to further popularize the use of liquid crystal displays, a reduction in cost must be achieved. Especially, it is required to reduce the cost of a color filter which occupies a large proportion of the total cost. Various methods have been tried to satisfy the required characteristics of color filters while meeting the above requirements. However, any method capable of satisfying all the requirements has not been established. The respective methods will be described below.\nThe first method is a pigment dispersion method. In this method, a pigment-dispersed photosensitive resin layer is formed on a substrate and patterned into a single-color pattern. This process is repeated three times to obtain R, G, and B color filter layers.\nThe second method is a dyeing method. In the dyeing method, a water-soluble polymer material as a dyeable material is applied onto a glass substrate, and the coating is patterned into a desired shape by a photolithographic process. The obtained pattern is dipped in a dye bath to obtain a colored pattern. This process is repeated three times to form R, G, and B color filter layers.\nThe third method is an electrodeposition method. In this method, a transparent electrode is patterned on a substrate, and the resultant structure is dipped in an electrodeposition coating fluid containing a pigment, a resin, an electrolyte, and the like to be colored in the first color by electrodeposition. This process is repeated three times to form R, G, and B color filter layers. Finally, these layers are calcined.\nThe fourth method is a print method. In this method, a pigment is dispersed in a thermosetting resin, a print operation is performed three times to form R, G, and B coatings separately, and the resins are thermoset, thereby forming colored layers. In either of the above methods, a protective layer is generally formed on the colored layers.\nThe point common to these methods is that the same process must be repeated three times to obtain layers colored in three colors, i.e., R, G, and B. This causes an increase in cost. In addition, as the number of processes increases, the yield decreases. In the electrodeposition method, limitations are imposed on pattern shapes which can be formed. For this reason, with the existing techniques, it is difficult to apply this method to TFTs. In the print method, a pattern with a fine pitch is difficult to form because of poor resolution and poor evenness.\nIn order to eliminate these drawbacks, methods of manufacturing color filters by an ink-jet system are disclosed in Japanese Patent Laid-Open Nos. 59-75205, 63-235901, and 1-217320. In these methods, inks containing coloring agents of three colors, i.e., R (red), G (green), and B (blue), are sprayed on a transparent substrate by an ink-jet system, and the respective inks are dried to form colored image portions. In such an ink-jet system, R, G, and B pixels can be formed at once, allowing great simplification of the manufacturing process and a great reduction in cost.\nWhen a color filter is to be manufactured by such an ink-jet system, an ink is discharged onto each pixel while an elongated ink-jet head having a plurality of ink discharging nozzles is scanned over a color filter substrate. This scanning is performed a plurality of number of times to color the respective pixel portions. In this case, the amounts of ink discharged from the respective ink discharging nozzles slightly differ from each other. If, therefore, each pixel array is colored with the same nozzle, the pixel arrays colored with the nozzles from which the ink is discharged in large amounts become dense in color, but the pixel arrays colored with the nozzles from which the ink is discharged in small amounts become light in color. Consequently, the resultant color filter has color irregularity.\nIn addition, to manufacture a color filter, such elongated ink-jet heads must be prepared for three colors, i.e., R (red), G (green), and B (blue). It takes a lot of time and labor to adjust the relative positions of these three heads."}
{"text": "1. Technical Field\nMethods and example implementations described herein are directed to interconnect architecture, and more specifically, automatically generating a Network on Chip (NoC) with pipelining solutions for the interconnects.\n2. Related Art\nThe number of components on a chip is rapidly growing due to increasing levels of integration, system complexity and shrinking transistor geometry. Complex System-on-Chips (SoCs) may involve a variety of components e.g., processor cores, DSPs, hardware accelerators, memory and I/O, while Chip Multi-Processors (CMPs) may involve a large number of homogenous processor cores, memory and I/O subsystems. In both SoC and CMP systems, the on-chip interconnect plays a role in providing high-performance communication between the various components. Due to scalability limitations of traditional buses and crossbar based interconnects, Network-on-Chip (NoC) has emerged as a paradigm to interconnect a large number of components on the chip. NoC is a global shared communication infrastructure made up of several routing nodes interconnected with each other using point-to-point physical links.\nMessages are injected by the source and are routed from the source node to the destination over multiple intermediate nodes and physical links. The destination node then ejects the message and provides the message to the destination. For the remainder of this application, the terms ‘components’, ‘blocks’, ‘hosts’ or ‘cores’ will be used interchangeably to refer to the various system components which are interconnected using a NoC. Terms ‘routers’ and ‘nodes’ will also be used interchangeably. Without loss of generalization, the system with multiple interconnected components will itself be referred to as a ‘multi-core system’.\nThere are several topologies in which the routers can connect to one another to create the system network. Bi-directional rings (as shown in FIG. 1(a)), 2-D (two dimensional) mesh (as shown in FIG. 1(b)) and 2-D Torus (as shown in FIG. 1(c)) are examples of topologies in the related art. Mesh and Torus can also be extended to 2.5-D (two and half dimensional) or 3-D (three dimensional) organizations. FIG. 1(d) shows a 3D mesh NoC, where there are three layers of 3×3 2D mesh NoC shown over each other. The NoC routers have up to two additional ports, one connecting to a router in the higher layer, and another connecting to a router in the lower layer. Router 111 in the middle layer of the example has both ports used, one connecting to the router at the top layer and another connecting to the router at the bottom layer. Routers 110 and 112 are at the bottom and top mesh layers respectively, therefore they have only the upper facing port 113 and the lower facing port 114 respectively connected.\nPackets are message transport units for intercommunication between various components. Routing involves identifying a path composed of a set of routers and physical links of the network over which packets are sent from a source to a destination. Components are connected to one or multiple ports of one or multiple routers; with each such port having a unique ID. Packets carry the destination's router and port ID for use by the intermediate routers to route the packet to the destination component.\nExamples of routing techniques include deterministic routing, which involves choosing the same path from A to B for every packet. This form of routing is independent from the state of the network and does not load balance across path diversities, which might exist in the underlying network. However, such deterministic routing may implemented in hardware, maintains packet ordering and may be rendered free of network level deadlocks. Shortest path routing may minimize the latency as such routing reduces the number of hops from the source to the destination. For this reason, the shortest path may also be the lowest power path for communication between the two components. Dimension-order routing is a form of deterministic shortest path routing in 2-D, 2.5-D, and 3-D mesh networks. In this routing scheme, messages are routed along each coordinates in a particular sequence until the message reaches the final destination. For example in a 3-D mesh network, one may first route along the X dimension until it reaches a router whose X-coordinate is equal to the X-coordinate of the destination router. Next, the message takes a turn and is routed in along Y dimension and finally takes another turn and moves along the Z dimension until the message reaches the final destination router. Dimension ordered routing may be minimal turn and shortest path routing.\nFIG. 2(a) pictorially illustrates an example of XY routing in a two dimensional mesh. More specifically, FIG. 2(a) illustrates XY routing from node ‘34’ to node ‘00’. In the example of FIG. 2(a), each component is connected to only one port of one router. A packet is first routed over the x-axis till the packet reaches node ‘04’ where the x-coordinate of the node is the same as the x-coordinate of the destination node. The packet is next routed over the y-axis until the packet reaches the destination node.\nIn heterogeneous mesh topology in which one or more routers or one or more links are absent, dimension order routing may not be feasible between certain source and destination nodes, and alternative paths may have to be taken. The alternative paths may not be shortest or minimum turn.\nSource routing and routing using tables are other routing options used in NoC. Adaptive routing can dynamically change the path taken between two points on the network based on the state of the network. This form of routing may be complex to analyze and implement.\nA NoC interconnect may contain multiple physical networks. Over each physical network, there may exist multiple virtual networks, wherein different message types are transmitted over different virtual networks. In this case, at each physical link or channel, there are multiple virtual channels; each virtual channel may have dedicated buffers at both end points. In any given clock cycle, only one virtual channel can transmit data on the physical channel.\nNoC interconnects may employ wormhole routing, wherein, a large message or packet is broken into small pieces known as flits (also referred to as flow control digits). The first flit is the header flit, which holds information about this packet's route and key message level info along with payload data and sets up the routing behavior for all subsequent flits associated with the message. Optionally, one or more body flits follows the head flit, containing the remaining payload of data. The final flit is the tail flit, which in addition to containing the last payload also performs some bookkeeping to close the connection for the message. In wormhole flow control, virtual channels are often implemented.\nThe physical channels are time sliced into a number of independent logical channels called virtual channels (VCs). VCs provide multiple independent paths to route packets, however they are time-multiplexed on the physical channels. A virtual channel holds the state needed to coordinate the handling of the flits of a packet over a channel. At a minimum, this state identifies the output channel of the current node for the next hop of the route and the state of the virtual channel (idle, waiting for resources, or active). The virtual channel may also include pointers to the flits of the packet that are buffered on the current node and the number of flit buffers available on the next node.\nThe term “wormhole” plays on the way messages are transmitted over the channels: the output port at the next router can be so short that received data can be translated in the head flit before the full message arrives. This allows the router to quickly set up the route upon arrival of the head flit and then opt out from the rest of the conversation. Since a message is transmitted flit by flit, the message may occupy several flit buffers along its path at different routers, creating a worm-like image.\nBased upon the traffic between various end points, and the routes and physical networks that are used for various messages, different physical channels of the NoC interconnect may experience different levels of load and congestion. The capacity of various physical channels of a NoC interconnect is determined by the width of the channel (number of physical wires) and the clock frequency at which it is operating. Various channels of the NoC may operate at different clock frequencies, and various channels may have different widths based on the bandwidth requirement at the channel. The bandwidth requirement at a channel is determined by the flows that traverse over the channel and their bandwidth values. Flows traversing over various NoC channels are affected by the routes taken by various flows. In a mesh or Torus NoC, there may exist multiple route paths of equal length or number of hops between any pair of source and destination nodes. For example, in FIG. 2(b), in addition to the standard XY route between nodes 34 and 00, there are additional routes available, such as YX route 203 or a multi-turn route 202 that makes more than one turn from source to destination.\nIn a NoC with statically allocated routes for various traffic slows, the load at various channels may be controlled by intelligently selecting the routes for various flows. When a large number of traffic flows and substantial path diversity is present, routes can be chosen such that the load on all NoC channels is balanced nearly uniformly, thus avoiding a single point of bottleneck. Once routed, the NoC channel widths can be determined based on the bandwidth demands of flows on the channels. Unfortunately, channel widths cannot be arbitrarily large due to physical hardware design restrictions, such as timing or wiring congestion. There may be a limit on the maximum channel width, thereby putting a limit on the maximum bandwidth of any single NoC channel.\nAdditionally, wider physical channels may not help in achieving higher bandwidth if messages are short. For example, if a packet is a single flit packet with a 64-bit width, then no matter how wide a channel is, the channel will only be able to carry 64 bits per cycle of data if all packets over the channel are similar. Thus, a channel width is also limited by the message size in the NoC. Due to these limitations on the maximum NoC channel width, a channel may not have enough bandwidth in spite of balancing the routes.\nTo address the above bandwidth concern, multiple parallel physical NoCs may be used. Each NoC may be called a layer, thus creating a multi-layer NoC architecture. Hosts inject a message on a NoC layer; the message is then routed to the destination on the NoC layer, where it is delivered from the NoC layer to the host. Thus, each layer operates more or less independently from each other, and interactions between layers may only occur during the injection and ejection times. FIG. 3(a) illustrates a two layer NoC. Here the two NoC layers are shown adjacent to each other on the left and right, with the hosts connected to the NoC replicated in both left and right diagrams. A host is connected to two routers in this example—a router in the first layer shown as R1, and a router is the second layer shown as R2. In this example, the multi-layer NoC is different from the 3D NoC, i.e. multiple layers are on a single silicon die and are used to meet the high bandwidth demands of the communication between hosts on the same silicon die. Messages do not go from one layer to another. For purposes of clarity, the present application will utilize such a horizontal left and right illustration for multi-layer NoC to differentiate from the 3D NoCs, which are illustrated by drawing the NoCs vertically over each other.\nIn FIG. 3(b), a host connected to a router from each layer, R1 and R2 respectively, is illustrated. Each router is connected to other routers in its layer using directional ports 301, and is connected to the host using injection and ejection ports 302. A bridge-logic 303 may sit between the host and the two NoC layers to determine the NoC layer for an outgoing message and sends the message from host to the NoC layer, and also perform the arbitration and multiplexing between incoming messages from the two NoC layers and delivers them to the host.\nIn a multi-layer NoC, the number of layers needed may depend upon a number of factors such as the aggregate bandwidth requirement of all traffic flows in the system, the routes that are used by various flows, message size distribution, maximum channel width, etc. Once the number of NoC layers in NoC interconnect is determined in a design, different messages and traffic flows may be routed over different NoC layers. Additionally, one may design NoC interconnects such that different layers have different topologies in number of routers, channels and connectivity. The channels in different layers may have different widths based on the flows that traverse over the channel and their bandwidth requirements.\nIn a NoC interconnect, if the traffic profile is not uniform and there is certain amount of heterogeneity (e.g., certain hosts talk to each other more frequently than the others), the interconnect performance may depend a lot on the NoC topology and where various hosts are placed in the topology with respect to each other and to what routers they are connected to. For example, if two hosts talk to each other frequently and need higher bandwidth, they should be placed next to each other. This will reduce the latency for this communication, and thereby reduce the global average latency, as well as reduce the number of router nodes and links over which the high bandwidth of this communication must be provisioned. Moving two hosts closer to one another may make certain other hosts far apart since all hosts must fit into the 2D planar NoC topology without overlapping with each other. Thus, tradeoffs must be made and the hosts must be placed after examining the pair-wise bandwidth and latency requirements between all hosts so that certain global cost and performance metrics is optimized. The cost and performance metrics can include the average structural latency between all communicating hosts in number of router hops, or the sum of the bandwidth between all pair of hosts and the distance between them in number of hops, or some combination thereof. This optimization problem is known to be non-deterministic polynomial-time hard (NP-hard) and heuristic based approaches are often used. The hosts in a system may vary is shape and sizes with respect to each other which puts additional complexity in placing them in a 2D planar NoC topology, packing them optimally leaving little whitespaces, and avoiding overlapping hosts.\nOne aspect of optimization of traffic profiles includes placement of pipeline stages as illustrated in FIG. 4. In a point to point example illustrated in FIG. 4, there are two routers R1 402 and R2 404 with input and output channels such as 406-1 and 406-2 between them. Each channel 406 can be configured to transmit data through 412-1 (from Router 1 to Router 2 using 406-1) and through 412-2 (from Router 2 to Router 1 using 406-2), which in turn is accompanied by a credit signal that flows in the opposition direction to the transmitted data (or control). Such credit signal is shown as 414-1 for channel 406-1 and 414-2 for channel 406-2. Each router 402/404 contains a clock, one or more buffers, and an internal flip flop for data management. Data transmitting router can further include a registering stage shown as 416-1 when data is transmitted from Router 1 to Router 2, and 416-2 when data is transmitted from Router 2 to Router 1. Furthermore, receiving router can include a FIFO shown as 418-1 when data is received by Router 2 from Router 1, and 418-2 when data is received by Router 1 from Router 2. To control timing or meet performance requirements, channels 406 between the routers 402/404 may employ one or more pipeline stages such as 408-1 and 408-2, collectively referred to as pipeline stages 408 hereinafter, within the channels 406. Pipeline stages 408 can include hardware elements such as flip-flops (e.g., JK, data/delay, etc.) to control the traffic flow. One or more pipeline stages 408 can be used as output registers such as 410-1 and 410-2, collectively referred to as output registers 410 hereinafter, which are employed at output channel of each router 402/404 to control traffic flow. A given router 402/404 may employ a single buffer for handling all traffic through the router, or can also employ a buffer for each input/output pair of channels managed by the router. Further, as the traffic flows through each of the pipeline stages 406, the traffic can be controlled to meet timing and or performance requirements based on flip-flop implemented at the pipeline stage. Pipeline stages 406 can further be configured to manage routing computation (RC), virtual channel allocation (VA), switch allocation (SA), switch traversal (ST) before it is delivered to the appropriate output port.\nIn the related art, there is no automated solution for placement of pipeline stages, and system designers may utilize the pipeline stages in a suboptimal manner when building a NoC. Further, as complexity of NoC increases, difficulty of determining placement and utilization of pipeline stages also increases."}
{"text": "1. Field of the Invention\nThis invention relates to photometric analysis of hydrogen sulfide in alkanol amines and, in particular, to a method and apparatus adapted to measure and to compensate for background photometric absorption caused, for example, by impurities in the amine stream so that an accurate measurement of the concentration of hydrogen sulfide may be obtained.\n2. Description of the Prior Art\nThe amine absorption process by which hydrogen sulfide is removed from refinery \"sour\" gas is well-known. The mechanism by which hydrogen sulfide is eliminated is understood and relates to the bonding of hydrogen sulfide with an alkanol amine. Typically, a \"lean\" amine, that is, an amine having an available site at which to bond with hydrogen sulfide, is introduced near the top of a countercurrent absorption tower, or absorber, while a sour gas stream is introduced near the bottom of the absorber. A chemical reaction within the absorber occurs wherein the hydrogen sulfide is bonded to the alkanol amine. The resultant bonded amine, or so-called \"rich\" amine, is then pumped to the inlet of a stripping column wherein the high temperature associated with high pressure steam is used to strip the hydrogen sulfide from the bonded amines. The stripped amines are conveyed, after cooling, to the absorption tower and the cycle is completed.\nThe amine-absorption desulfurization process is described in a paper written by R. S. Saltzman entitled \"Continuous Photometric Analysis for Measurement and Control in Sulfur Recovery Operations\" presented at the Texas A&M University Instrumentation Symposium for Process Industries, January 1977.\nFor various economic reasons the amount and cost of energy utilized to produce high pressure steam used in the stripping column has become the focus of conservation activity. When energy costs and availability were less important considerations it was the typical practice to use steam to lower the concentration of hydrogen sulfide to a level below that necessary for efficient absorption in the absorber. However, as the cost of energy has increased, this practice is no longer economical. It is, accordingly, important to provide some indication as to the level of hydrogen sulfide remaining in the effluent from the stripping column so that the efficiency of the stripping column may be optimized from an energy consumption standpoint.\nTo this end it is known that photometric analyzers have been utilized to provide an indication of the amount of hydrogen sulfide remaining in the effluent of the stripping column. One such apparatus is described in the above-mentioned article and is sold by E. I. du Pont de Nemours and Company as the 400 Photometric Analyzer. Such apparatus relies upon the strong ultraviolet radiation absorption characteristic of the alkanol amine-hydrogen sulfide bond to measure the hydrogen sulfide concentration in the stripping column effluent. In general, two problems have been observed with regard to obtaining accurate measurements of the hydrogen sulfide concentration.\nThe first problem relates to the acute temperature sensitivity of the absorption characteristic of the amine-hydrogen sulfide bond. At higher temperatures, more ultraviolet radiation absorption occurs. Accordingly, it has been necessary to provide a stabilizing heat exchanger in the monitoring path upstream of the photometric analyzer to insure that the sample of the effluent from the stripping column introduced into the sample cell of the analyzer is maintained at a predetermined constant temperature. By this expedient it has been found that the effect of temperature on the ultraviolet radiation absorption characteristic of the amine-hydrogen sulfide bond is minimized.\nA further problem relates to the effects of recycled impurities in the amine loop. As the amines become recycled repeatedly it is possible that dirt, corrosion inhibitors, or other contaminants produced by chemical breakdown may occur to form such impurities. Such impurities behave as absorbers of ultraviolet radiation at the same wavelengths as are absorbed by an amine-hydrogen sulfide bond. Accordingly, incorrect measurements of the concentration of hydrogen sulfide in the effluent from the stripping column can occur. This could result in the unnecessary increase in steam use with a concomitant increase in energy consumption. Increased energy consumption under these circumstances would not be economically justified and would, therefore, be disadvantageous. It should be noted that the application of the device mentioned in the referenced article was on specially controlled processes where impurities in the stream were maintained at low levels and the sample temperature was closely regulated.\nTo obtain a measurement of the hydrogen sulfide concentration in a stream it is known also to introduce a sample of the stream that contains bonded hydrogen sulfide-amine into a vessel and to thereafter introduce a predetermined amount of hydrogen chloride into the vessel. The hydrogen chloride replaces the hydrogen sulfide and the hydrogen sulfide is measured to provide an indication of the concentration of hydrogen sulfide in the sample."}
{"text": "It is known in the art that the presence of biogenic amines in human body fluids may reveal or suggest pathological conditions and dysfunctions. For example, elevated levels of certain biogenic amines in urine may indicate the presence or the likelihood of the presence of a cancer (there are many papers dealing with this—see, for instance, Suh, J W, Lee, S H, Chung, B C, Park, J, Urinary Polyamine Evaluation for Effective Diagnosis of Various Cancers, Journal of Chromatography B, 1997, Vol. 688, Iss 2, pp. 179–186). Several of the types of vaginal diseases may be expressed in elevated levels of biogenic amines in vaginal discharge and fluids (see, for instance, C. S. Chen, R. Amsel, D. A. Eschenbach and K. K. Holmes, Biochemical diagnosis of vaginitis: determination of diamines in vaginal fluid, J. Infectious Disease 145 (1982), pp. 337–345).\nBody fluids may include e.g. urine, blood, serum, saliva, vaginal discharge and fluids, etc. Further, samples in which the presence of biogenic amines may be determined may not be fluids, but, e.g., skin and tissues, swipe samples, etc. Even direct sniffing of skin or breath exhaled by a subject may provide information in this respect. This should be understood whenever body fluids are mentioned in this application.\nChemical changes in the living system or degradation processes of cells after death are accompanied with formation of molecular byproducts. These processes include the breaking down of peptides and DNA strands to smaller components, and changes in the building blocks, amino acids, leading to the formation of amines. Not only amino compounds are produced, but other, smaller molecules, such as aldehydes and alcohols are also formed. One of the processes of particular interest is the breakdown of amino acids and the production of diamines and polyamines. For example, decarboxylation of histidine, ornithine, lysine, produces histamine, putrescine and cadaverine respectively.\nSeveral analytical methods have been proposed for the analysis of biogenic amines. Most of these are laboratory methods that require expensive equipment, extensive sample preparation and the skills of a trained analytical chemist or technician. Among these are high performance liquid chromatography (HPLC), or gas chromatography after derivatization of the samples. Biosensors may also be used, as well as various spectrometric techniques. Solid state sensors have also been proposed, but generally lack specificity.\nIon Mobility Spectrometry (also, briefly, IMS) is a known analytical method and its application for the determination of aliphatic and aromatic amines has been suggested: see, for instance, Z. Karpas, Ion Mobility Spectrometry of Aliphatic and Aromatic Amines, Anal. Chem. 61 (1989), 684. An apparatus for carrying out this method—the Ion Mobility Spectrometer (IMS)—is used primarily for detection, identification and monitoring of trace amounts of gases and vapors. It is particularly suitable for detection of compounds that have high proton affinity and form stable positive ions, or for compounds that have a high electronegativity and readily form stable negative ions. IMS is fully discussed in J. I. Baumbach and G. A. Eiceman, Appl. Spectrosc. 1999, vol.53, pp.338A–355A. However, any device that may be used for determining or measuring the mobility of ions may be used for carrying out the invention, and therefore any reference to IMS in this description and claims should not be construed as a limitation, but should be construed any including instrument for determining or measuring the mobility of ions.\nThe knowledge of the prior art as to the importance of biogenic amines for the possible detection of pathological conditions and as to the analysis of biogenic amines, including the use of IMS, has failed so far to provide a simple and reliable method for the diagnosis of vaginal disorders, particularly, though not exclusively, bacterial vaginosis that affects a large number of women. The provision of such a diagnostic method would constitute a valuable contribution to the medical art. However, such a method is not available: the detection of biogenic amines in vaginal fluids is known to suggest the presence of a pathological condition, but it does not provide specific and reliable information and merely suggests to the specialized physician the desirability of carrying out whatever tests and examinations may finally lead to a diagnosis.\nFurther, the present knowledge does not provide the practicing physician with an apparatus for the quick diagnosis of bacterial vaginosis and other pathological conditions, by a simple and direct way and without the application of knowledge and technology that are typical of different branches of science and are not found together in any physician, no matter how competent and dedicated. It would be extremely valuable to provide an instrument and method whereby the average physician could obtain from bodily fluids, quickly and in a reliable way, a diagnostic indication of specific diseases and/or pathological conditions, even though such a diagnostic indication may not be final and conclusive and may require, whether positive or negative, verification and integration.\nIt is therefore a purpose of this invention to provide method for the diagnosis of vaginal disorders, particularly, though not exclusively, bacterial vaginosis.\nIt is another purpose to provide a method for carrying out such diagnosis automatically and in real time.\nIt is a further purpose to provide, automatically and in real time, information of fundamental value in the diagnosis of a variety of pathological conditions.\nIt is a still further purpose to provide an apparatus for the quick diagnosis of bacterial vaginosis and other pathological conditions.\nIt is a still further purpose to provide a portable instrument that is capable of carrying out the spectrometry of bodily fluids and automatically derive from said spectrometry significant diagnostic indications.\nIt is a still further purpose to provide such an instrument that can be directed to provide diagnostic indications for specific diseases and/or pathological conditions.\nIt is a still further purpose to provide such an instrument that consists of the combination of components known in the art and readily available.\nIt is a still further purpose to provide such an instrument that can be widely used by physicians and medical institutions and is not excessively expensive.\nIt is a still further purpose to provide such an instrument which can be used, with the due precautions and warnings, by persons other than physicians and even by the patients themselves.\nIt is a still further purpose to provide such an instrument which can be used for purposes that are not diagnostic purposes, but are relevant to the public health, for instance, the control of the condition of food, such as, but not exclusively, the freshness of meat, fish and their products, as well as and seafood.\nIt is a still further purpose to provide such an instrument that permits instant examination of tissues removed during an operation as an indication for malignant tissues."}
{"text": "Many types of electronic devices are battery operated and thus preferably consume as little power as possible. An example is a cellular telephone. Further, it may be desirable to implement various types of multimedia functionality in an electronic device such as a cell phone. Examples of multimedia functionality may include, without limitation, games, audio decoders, digital cameras, etc. It is thus desirable to implement such functionality in an electronic device in a way that, all else being equal, is fast, consumes as little power as possible and is as efficient as possible. Improvements in this area are desirable."}
{"text": "High voltage equipment is commonly sealed within a welded enclosure. The seal is required because the enclosure contains either a gaseous (e.g., SF.sub.6) or liquid (e.g., oil) dielectric atmosphere within which the equipment operates. To enable connection between the enclosed equipment and external power lines, a bushing is used to provide a conductor through the wall of the enclosure while simultaneously insulating the conductor from the enclosure and maintaining the seal.\nBushings typically consist of a unitary metal conductor encapsulated in epoxy insulation. The inside end of the bushing (i.e., the proximal end of the conductor) is usually connected rigidly (e.g., by a metallic clamp and hardware) to the enclosed high voltage equipment, while the outside end of the bushing is designed to matingly receive a connector device, such as an elbow connector. The outside end of the bushing (i.e., the insulator body and/or the bushing-to-elbow interface) is subject to damage during both transportation and use which quite often renders the bushing unusable. In this event, the bushing must be replaced, which for a unitary bushing requires opening of the sealed enclosure and disconnection of the bushing from the high voltage equipment. Such replacement procedures are costly, time-consuming and often quite difficult to perform, especially when the equipment is located in the close quarters of an underground vault."}
{"text": "The present invention relates to heating systems, and more specifically to a heat recovery system for boiler and domestic water.\nConventionally, the cold water supply line is fed to a boiler tankless coil where the cold domestic water is heated by the boiler water and released to be used in hot water fixtures of a house. The boiler water, on the other hand, when returned from the house heating units or house radiating system, returns to the boiler through a water circulator. After being reheated, the boiler water is forced back into the house heating units.\nA major disadvantage of the conventional arrangement is that feeding cold water to the boiler tankless coil reduces the efficiency of the boiler since the boiler must heat the cold water from approximately 40.degree.-50.degree. F. to approximately 120.degree.-180.degree. F. Although requiring a great input of heat for this purpose, a great deal of the boiler heat escapes through the stack with the flue gases without the heat being retrieved.\nIn addition to reducing the efficiency of the boiler for heating the boiler water, prior art systems which only utilized an internal boiler tankless coil were also deficient in providing sufficient quantities of domestic water for bathing, using major appliance, etc. It frequently occurs that the capacity of the boiler tankless coil is insufficient for producing sufficient quantities of domestic hot water or heating such water sufficiently rapidly to provide the desired quantities.\nIn U.S. Pat. No. 2,189,749, which discloses a water heater, suggests the use of a supplemental heating unit using flue gases passing through an envelope filled with boiler water through which a coil for domestic water passes. However, the arrangement described requires that the flue gases be forced from the heating boiler down into and through the supplemental heating unit. Such an arrangement is not recommended for many reasons. Firstly, it causes back pressure on the system as heated gases rise to be vented. Also, the excess pressure can cause relief-type devices to blow off which can soot up an entire dwelling. Further, with the arrangement in the patent, using a tube and baffle design, soot will collect if the oil burner is not running to peak performance. Additionally, the disclosed device utilizes a damper operated by an actuating rod, and the damper could carbonize and that sends heat causing the butterfly to malfunction. This could also cause back pressure which may result in a release of soot in the premises.\nU.S. Pat. No. 2,554,338 is for a water heater and uses an auxiliary water heating tank inside the furnace with a pipe coaxial with the flue leading to a tank surrounding the flue. However, the two auxiliary hot water tanks which are disclosed are disposed directly in coal which is transferring heat to the domestic water. Using oil as a fuel and placing the auxiliary hot water tanks in the flames of a oil fired boiler would cause many problems, form soot, smoke and in most cases, could not fit into conventional boilers used today. Furthermore, to install domestic water lines through the smoke pipe and have gases pass through a domestic water tank is not practical as it gravitates to the tank and must be heated to move the water. For these reasons, the water heater disclosed in this patent could not be applied to today's conventional oil heating units.\nA furnace system for heating air and water is disclosed in U.S. Pat. No. 2,827,893. This system provides a water heater in the exhaust flue in conjunction with a circulation system and a forced air heating system. A summer and winter changeover is used to eliminate the use of the separate water heater in winter and using cold water in the coils in the summer to aid cooling of a building. However, the use of a coil to extract heat from a hot air unit or heat exchanger is very inefficient since the heat is being taken from the home heating air. The use of a separate circulator disclosed in this patent to move water is an additional problem for those who service the unit. Further, the coil is restrictive both to the flow of hot air which is rising and the return air. The patent also discloses the use of a storage tank and expansion tank. However, the use of these is not practical since there are too many parts to be installed. The use of a unit of this type, which is complex, along with a hot water heater is not at all practical as it offers no saving utilizing such a system. The disclosed system merely retards heat flow in the unit and takes heat from it as well.\nAn energy recovery and storage system is also disclosed in U.S. Pat. No. 4,037,786. This patent teaches an arrangement which provides a heat recovery and storage system utilizing heat transfer pipes and circulating water from a storage tank or heat sink in response to sensors responding to water temperature and flue gas temperature. However, the use of transfer heating rods to heat domestic water is inefficient as the heat only strikes the surface area of those rods. Additionally, the use of a hot water heater and storage tank is costly for the operation of the hot water heater and the extra boost given to it by the energy recovery and storage system is grossly inefficient. The use of a storage tank along with the hot water heater is not practical for the space needed to install such a unit.\nIn U.S. Pat. No. 3,896,992 there is disclosed a heat recovery system for space heating and for potable water heating. This system provides a heat recovery coil and a flue that may alternately provide additional heat to the heating system or provide preheating of the domestic hot water system. However, the coil inside the flue pipe will restrict flue gas and reduce draft causing soot conditions. Also, dismantling the flue pipe for cleaning can create a substantial problem. Further, the water pre-heater tank along with the disclosed pressure gauge and pressure relief valve all require additional space. The disclosed system is extremely elaborate, too costly and impractical."}
{"text": "This invention relates to a method of evaluating data pertaining to respiration and heart beat and an apparatus for evaluating them, especially for forecasting sudden infant death syndrome (SIDS) by investigating the correlation between respiration and heart beat in a normal state and respiratory standstill in sleep.\nGenerally, the data from organisms includes EEG, cardiac electricity, respiration, ocular movement, and EMG. Signal forms of this data are relatively more useful than from EEG etc. It is therefore easier to analyze an R--R interval as a change of peak interval of an ECG wave, chronologically.\nConsidering that respiratory arrhythmia depends upon an efferent impulse from an extension receptor, standstill of the efferent impulse from an extension receptor in respiratory arrest in sleep will have an effect on heart beat. It has been already reported that in Cheyne Stokes respiration during hyperrespiration the increase of heart beat and decrease of blood pressure are observed; during its apnea the decrease of heart beat and increase of blood pressure occur.\nHere, sleep-apnea syndrome is defined as three following cases:\n(1) At least 30 episodes of apnea appear in the night sleep of 7 hours.\n(2) An apnea episode appears in REM-sleep as well as in non-REM-sleep.\n(3) An episode of apnea has a tendency to appear continuously.\nThe classification of sleep apnea syndrome according to Gastant is divided into the three following types\n(1) Central type: A respiratory standstill at the nostrils and mouth appears along with the disappearance of thoracid and abdominal movements. This indicates the stoppage of activity of the respiratory center.\n(2) Obstructive or upper tracheal type: In spite of respiratory efforts by thracic and abdominal movements, respiration stops, suggesting an obstruction in the upper tracheal tract.\n(3) Mixed type: An intermediate type between central and obstructive ones. Although at the first phase of apnea the thracic and abdominal movements cease, its symptom becomes more severe gradually.\nIn the central type, the apnea episode often appears repeatedly in an almost constant period, frequently observed in non-REM-sleep, but in the obstructive type an apnea episode appears at various intervals and its duration is also not constant. Also, in the central type, during hyperventilation in Cheyne Stokes respiration an increase of heart beat is observed. During the apnea phase its decrease and a periodical change are noticed.\nThe mechanism of the central type of sleep-apnea can be explained by neurogenous or respiratory disturbance theory. The neurogeneous theory is based on the assumption that there is an oscillation center in the brainstem, and an inhibitory mechanism for controlling oscillation in the forebrain. The respiratory disturbance theory rests on the basis that a respiratory disturbance delays a feed-back loop between the respiratory center and chemical receptor, and an oscillation phenomenon appears.\nIt has not been determined which effect really works, or whether both effects work simultaneously."}
{"text": "1. Field of the Invention\nThe present invention relates to a head suspension, a method of manufacturing a base plate of the head suspension, and a method of manufacturing the head suspension, the head suspension being a device to support a read/write head to write and read data to and from a magnetic disk or a hard disk in a disk drive of for example, a personal computer.\n2. Description of Related Art\nA base plate for a head suspension is manufactured as illustrated in FIG. 12, for example. FIG. 12 illustrates a method of manufacturing a head suspension according to a related art (Japanese Unexamined Patent Application Publication No. 2012-14812). In FIG. 12, a base plate chain 101 is processed into a head suspension chain 103, and from the head suspension chain 103, individual head suspensions 105 are cut off.\nThe base plate chain 101 includes base plates 107 each of which is fabricated into the head suspension 105. The head suspension 105 is of a dual actuator system that uses, in addition to a voice coil motor to turn the head suspension 105 as a whole in a sway direction or a width direction, a piezoelectric element to minutely turn a head of the head suspension 105 in the sway direction. The base plate 107 includes a boss 111 and an actuator base 113. In the base plate chain 101, a plurality of the base plates 107 are joined to a frame 112.\nThe actuator base 113 has a rear end 115 and a front end 117. The rear and front ends 115 and 117 are spaced from each other in a direction of a rotation radius of the head suspension 105 and are connected to each other through a link 119 that is present at a widthwise center of the base plate 107. On each side of the link 119, there are a pair of actuator attaching areas 121a and 121b that are openings in each of which the piezoelectric element is attached.\nIn this specification, a “longitudinal direction” corresponds to the direction of the turning radius of the head suspension and a “width direction” to a direction orthogonal to the longitudinal direction.\nAlong each side of the rear end 115, there are side segments 123b and 125b, and along each side of the front end 117, there are side segments 123a and 125a. The side segments 123a (125a) and 123b (125b) extend in the longitudinal direction, are spaced from each other, and face each other, so that the front end 117 is movable relative to the rear end 115.\nThe base plate chain 101 is deburred by, for example, chemical polishing, electrolytic polishing, or barrel polishing, and load beams 125 are connected to the respective base plates 107 in the base plate chain 101. Thereafter, piezoelectric elements 127a and 127b each made of, for example, PZT (lead zirconate titanate) are attached in the respective actuator attaching areas 121a and 121b of each base plate 107 with a nonconductive adhesive. Thereafter, a flexure having a head and wiring is attached to each head suspension 105, thereby forming the head suspension chain 103. The head suspension chain 103 is cut into individual head suspensions 105.\nIn the base plate chain 101 to be subjected to the above-mentioned processes, each base plate 107 has the actuator attaching areas 121a and 121b and separated side segments 123a, 123b, 125a, and 125b. These parts may entangle among a plurality of base plate chains 101 during the transportation of the base plate chains 101 between processes and may be deformed by receiving external force. The deformation causes material and product losses and increase costs. To avoid the entanglement and deformation, the base plate chains 101 may be aligned before deburring or heat treatment. This, however, deteriorates productivity.\nIf the base plates 107 are thinner to reduce their weight, the above-mentioned problem will become more serious.\nFIG. 13A is a plan view illustrating a base plate 107A with outer frames according to the related art and FIG. 13B is a plan view illustrating a base plate 107B without outer frames according to the related art. These base plates 107A and 107B are individually manufactured. Namely, they are not chained.\nIn FIG. 13A, the base plate 107A has a single actuator attaching area 121 being an opening and outer frames 129a and 129b that outwardly protrude with respect to the actuator attaching area 121 in the width direction. The actuator attaching area 121 receives a single piezoelectric element.\nIn FIG. 13B, the base plate 107B is similar to the base plate 107 of FIG. 12, and therefore, is illustrated with the same reference numerals as those of FIG. 12.\nWhen a plurality of the base plates 107A (107B) are simultaneously processed through deburring, heat treatment, and head suspension assembling, the actuator attaching areas 121 (121a, 121b) may entangle among the base plates 107A (107B), to deform front ends 117A (117) and the like.\nThe barrel polishing, in particular, is speedier for the deburring of the base plates 107A (107B). It, however, frequently causes the entanglement and deformation of the base plates 107A (107).\nCompared with the base plate 107A of FIG. 13A, the base plate 107B of FIG. 13B has no outer frames, and therefore, has low rigidity and easily entangles with others to increase material and product losses and decrease productivity.\nTo solve the problem, the base plate 107B without frames may be provided with a scrap section in a half-finished stage.\nFIG. 14 is a view illustrating Step S101 to form a half-finished base plate 107BA and Step S102 to form the base plate 107B.\nStep S101 of FIG. 14 forms the half-finished base plate 107BA by pressing a stainless steel thin plate of 100 to 250 micrometers thick. The half-finished base plate 107BA has scrap sections 131a and 131b which are black-colored in FIG. 14 for the sake of explanation. In practice, the scrap sections 131a and 131b are made of the same material as the half-finished base plate 107BA and are integral therewith with no boundaries.\nA plurality of the half-finished base plates 107BA are collectively deburred by, for example, barrel polishing, are heat-treated, and are subjected to Step S102. The deburring and heat treatment are carried out to the half-finished base plates 107BA each having the scrap sections 131a and 131b that increase the rigidity of the actuator attaching areas 121a and 121b. Even if many half-finished base plates 107BA are simultaneously processed by, for example, barrel polishing, the actuator attaching areas 121a and 121b of each half-finished base plate 107BA will not entangle with others.\nStep S102 of FIG. 14 cuts off the scrap sections 131a and 131b by pressing from the half-finished base plate 107BA, to form the base plate 107B that integrally has the actuator attaching areas 121a and 121b and the boss 111.\nThe related art of FIG. 14 is capable of suppressing the deformation of the half-finished base plate 107BA during the deburring, heat treatment and the like and realizing easy handling of the half-finished base plate 107BA for the deburring, heat treatment and the like.\nThe related art, however, has a problem that Step S102 carried out after the deburring process of Step S101 leaves burrs or cut burrs on cutting surfaces at which the scrap sections 131a and 131b are cut off from the half-finished base plate 107BA. These cut burrs may drop off in a completed hard disk drive into which a head suspension having the base plate 107B is installed.\nIf many half-finished base plates 107BA are processed with a metal mold for cutting off the scrap sections 131a and 131b, blades of the metal mold will become dull to enlarge the cut burrs. The large cut burrs are likely to drop off in hard disk drives.\nTo minimize the cut burrs, the blades of the metal mold must frequently be maintained. This, however, deteriorates productivity due to intermittence of the manufacturing process."}
{"text": "An optical collimator is a device which collects rays from a point source of light such as a light bulb or light emitting diode and causes those rays to emerge in parallel from a surface. Examples of collimators include lenses or curved mirrors found in a flashlight or car headlamp. In these examples, a volume of space exists between the point source and the surface from which the collimated light exits. In some use environments, this space may be inconvenient, as it may increase the overall size of an optical device that utilizes the collimator."}
{"text": "On removing oils and the like such as machine oil, rolling oil, and rust-preventive oil attached to machines, mechanical parts, and the like, or on removing oils and the like generated during fiber-processing stage and attached to facilities, there are used in recent years a water-base cleaner consisting mainly of nonionic and anionic surfactant, and a strong-alkali cleaner combining the water-base cleaner with a strong alkali. In case of conventional cleaners which enhance the increased cleaning power and foamability, however, the water-oil separability in the cleaned waste liquid becomes insufficient to apply heavy load to the wastewater treatment, because oils form stable emulsion in the cleaned waste liquid after cleaning or because oils stably become solubilized in the cleaned waste liquid after cleaning. Current serious concern about the environmental problems raises the requirement of reducing the quantity of waste liquid generated from the water-oil separation treatment and of recycling the cleaner. Thus a cleaner having excellent water-oil separability in the cleaned waste liquid is desired.\nAs the above cleaners, there is disclosed a cleaner composition containing (A) a polyoxyalkylene additive of alkyl amine or alkyl amide, having hydrocarbon group with the number of carbon atoms from 8 to 22, and (B) at least one compound selected from the group consisting of an alkyldimethylamine oxide having hydrocarbon group with the number of carbon atoms from 8 to 22, an alkyl-di (aminomethyl)glycine having hydrocarbon with the number of carbon atoms from 6 to 22, a straight chain sodium dodecylbenzene sulfonate, and betaine lauryldimethyl aminoacetate. Furthermore, there is a disclosed cleaner prepared by mixing the (A) component with the (B) component at a specific mixing rate intending to improve the cleanability, the foamability, and the water-oil separability after cleaning, (for example, see Patent document 1.)\nIn addition, there is a disclosure of water-oil separating cleaner containing an N-alkyl-N-(2-hydroxyalkyl)-iminoethylene carbonate, (for example, see Patent document 2.)    [Patent document] Japanese Patent Laid-Open No. 2003-119496    [Patent document] Japanese Patent No. 3430147\nRegarding the cleaning composition disclosed in Patent. document 1, however, a cleaner composition which uses a polyoxyethylene additive of alkylamine or alkylamide having a polyoxyethylene chain with oxyethylene units from 2 to 50, particularly preferably oxyethylene units from 10 to 50, as the (A) component becomes emulsion by itself to show water-soluble property. Thus, when that type of composition is used to the machines, the mechanical parts, and the like to which oils such as grease, engine oil, machine oil, and rust-preventive oil, (in particular grease-base oil, naphthene-base mineral oil, paraffin-base mineral oil, poly-α-olefin, polyol ester, and polydimethyl siloxane), are attached, the satisfactory cleanability and the water-oil separability in the cleaned waste liquid cannot be attained at the same time, and further a problem of rust-generation on the parts to be cleaned arises during or after cleaning.\nAs for the use of cleaner disclosed in Patent document 2, the water-oil separability after cleaning is significantly influenced by the temperature of the cleaning liquid. Since, at low temperatures, separation of oil layer from aqueous one takes a long time, the reduction in load to the wastewater treatment is not sufficient. In addition, cleanability and rust-preventive performance are not sufficient."}
{"text": "A “universal” prior art motor, such as is shown in FIG. 1, is an electric series wound motor which can run on AC or DC power. Such motors are small, frequently connected directly to standard 110 volt AC power and run at high speed. They are generally employed in hand-held and counter top appliances such as mixers, hand drills, vacuum cleaners and in gardening tools such as edge trimmers and leaf blowers.\nSuch motors use two field coils wound around two poles for creating a magnetic field in the stator and a rotor with many coils (usually twelve) “wave” wound around the rotor in the rotor slots so that parts of the rotor coil conductor turns lie parallel to the rotor axis. The rotor coils are connected to the commutator pads (usually twelve). Two brushes, which are connected in series with the stator field coils, contact the commutator. The power is connected to the field coils and the speed of the motor is varied by a simple power supply which varies the voltage. When an electric current flows through the conductors of the rotor coils, which are in the magnetic field, the conductors feel a force which turns the rotor about its axis.\nNo other design has yet as effectively satisfied the special niche need for high speed, variable speed and low cost in small appliances as has the “universal” motor. Typically the speed required in such devices is above 10,000 rpm. Such high speed electric motors deliver high power in a relatively small lightweight package at a low price. The “universal” motor also meets the requirement for smaller motors in hair blowers. In these devices, the motor usually has permanent magnets replacing the stator field coils.\nRegardless of whether a stator has coils or permanent magnets for generating magnetic flux, the coils in the rotor also produce a magnetic field which according to some literature “distorts and opposes” the main stator field. This is termed “rotor reaction” and is considered a detriment to efficient motor action.\nIn the currently used universal motors, with field coils, the longitudinal length of the conductors in the wound rotor is an important parameter because it is the interaction of the conductors with the stator's magnetic field which produces the torque. This torque is found to be approximately proportional to the longitudinal length of the rotor conductors. Therefore, the “bank” of coils in the wound rotor need only be a single “bank” and be whatever length produces the torque desired. The wound rotor purpose is to provide as many conductors, arranged longitudinally, as can be fitted in the space. If there were some other way to do this other than making coils it may still function. The coils are incidental, not essential.\nCurrently, available universal motors use brushes which are linearly guided and must be connected to electrical power via flexible pigtails. This requires an arrangement which is costly, fragile and subject to breakage and is difficult to service.\nTwo other types of electric motor have been built. They are the variable reluctance, and the switched reluctance motors. These motors have advantages of low cost through the elimination of unnecessary copper but have severe disadvantages of torque variation (ripple) and acoustic noise which, so far, is only partly mitigated by expensive tailored electronically controlled drives."}
{"text": "1. Field of the Invention\nThe invention is in the field of computing systems and more specifically in the field of network security.\n2. Description of the Related Art\nSeveral approaches to controlling external access to computing networks have been developed. The goals in developing these approaches include limiting access to authorized users and assuring that computing devices employed by these users do not include malicious computing code such as viruses, worms, or Trojan horses. The need for access control has grown with users' demands for accessing secure networks over the Internet and from personal devices such as laptop computers and personal digital assistants.\nA first level of access control is achieved by requiring authentication of a user. This may be accomplished by requiring the user to enter a username and password or by reading a MAC address or other identifying information from an access device. In some systems, a network switch is programmed to grant access to a secure network only after proper authentication is achieved. Systems capable of using this approach include those using a proprietary VMPS protocol from Cisco Systems, Inc. of San Jose, Calif. or using a IEEE 802.1x standard protocol.\nReliance on mere user authentication includes several disadvantages. For example, there are no provisions to assure that an access device used by the authenticated user meets network security policies. The access device may have out-of-date virus software, may have security vulnerabilities, or may be otherwise compromised. Further, this approach requires that access points (e.g., network switches) support one of a specific set of access protocols. Because many access points do not support VMPS, 802.1x, or a similar protocol, implementation of this approach on a large preexisting network, such as a corporate network, may be prohibitively expensive.\nA greater level of access control may be achieved by including a gatekeeper between the secure network and the access device. The gatekeeper is configured to ensure that the access device conforms to a predetermined security policy. FIG. 1 is a block diagram of a Secure Network, generally designated 100, including two instances of a Gatekeeper 110 each associated with a different instance of an Access Point 120. Secure Network 100 further includes, as an illustrative example, a network File Server 130, a Network PC (personal computer) 140 and a Printer 150 included in and connected via, a local area network designated LAN 125.\nEach instance of Access Point 120 may include a large number of individual communication ports. These ports are used to connect through GateKeeper 110 to other elements within Secure Network 100 such as File Server 130 or LAN 125. Large corporate networks may include many instances of Access Point 120, each including hundreds of individual communication ports and being associated with an instance of GateKeeper 110.\nSome instances of Access Point 120 are capable of supporting virtual local area networks (VLANs). VLANs are generated by associating communication ports within Access Point 120 with separate virtual networks. For example, one Access Device 160 may be placed on a different VLAN than another Access Device 160 by assigning the communication ports of Access Point 120 to separate VLANs. From the point of view of these devices, the effect of a VLAN is equivalent to having a separate hardwired network.\nGateKeeper 110 is configured to control access to Secure Network 100. In addition to authenticating users who wish to access Secure Network 100 using an Access Device 160, GateKeeper 110 is configured to ensure that Access Device 160 conforms to a predetermined security policy, before granting access to Secure Network 100. For example, GateKeeper 110 may make certain that Access Device 160 has up-to-date virus software and encryption protocols as proscribed by the security policy. Once GateKeeper 110 has verified that Access Device 160 satisfies the security policy, Access Device 160 is allowed to communicate through GateKeeper 110 to LAN 125.\nA disadvantage of the use of GateKeeper 110, as practiced in the prior art, is that all communications between Access Device 160 and LAN 125 pass through GateKeeper 110. For large networks including many instances of Access Point 120, each of which may include many communication ports, this can be a significant burden. The use of one instance of GateKeeper 110 to support numerous instances of Access Device 160 is limited by bandwidth and required sophistication (e.g., cost) of GateKeeper 110. It is undesirable for GateKeeper 110 to become a limiting factor on the bandwidth of communication between instances of Access Device 160 and Secure Network 100. Further, the use of a separate GateKeeper 110 for each instance of Access Device 160 or Access Point 120 is often prohibitively expensive. Even if a separate GateKeeper is used for each Access Point 120, there is no isolation between compliant and non-compliant Access Devices 160 which are attached to the same Access Point 120.\nThere is, therefore, a need for improved systems and methods of controlling access to secure networks."}
{"text": "1. Field of the Invention\nThe present invention relates to a chemically amplified resist composition, and more particularly, to a photosensitive polymer having a cyclic backbone, and a resist composition for an ArF excimer laser obtained therefrom.\n2. Description of the Related Art\nAs semiconductor devices become highly integrated and complicated to fabricate, fine pattern formation is required.\nFurther, as the capacity of a semiconductor device increases to exceed 1 giga bit, a pattern size having a design rule of less than 0.2, .mu.m is required. Accordingly, there are limitations in using a conventional resist material with a KrF excimer laser (248 nm). Thus, a new resist material capable of being developed using an ArF excimer laser (193 nm) has been developed in a lithography process.\nThe resist material used in the lithography process using the ArF excimer laser has several problems in being commercially used, compared to the conventional resist materials. The most typical problems are transmittance of a polymer and resistance to dry etching. As the widely known ArF resist materials, (meth)acrylate polymers are generally used. In particular, the most typical resist material is a poly(methyl methacrylate-tert-butyl methacrylate-methacrylic acid) terpolymer system manufactured by IBM, Inc. However, such polymers have very weak resistance to dry etching.\nAccordingly, to increase the resistance to dry etching, a polymer having a backbone composed of an alicyclic compound such as an isobornyl group, an adamantyl group or a tricyclodecanyl group, is used. However, the resulting resist still exhibits weak resistance to dry etching."}
{"text": "Windows that can transmit a controllable fraction of incident light intensity would be beneficial in a variety of applications. Thus, windows of buildings or vehicles, windows between rooms, and personal items such as glasses or goggles are often used in situations in which it would be advantageous if their optical transparency could be adjusted, for example by electrical means.\nGreat efforts have been expended to improve the ways and means of selectively controlling the transmission of light through window structures. A common approach to light control involves using an opaque window shade to reduce the transmission of electromagnetic radiation. Such shades may either be purely mechanical (the most common type) or may be controlled by a motor. Another approach to variable control of light transmission can be achieved by mechanically rotating a pair of polarizing films where the relative angle between polarizing axes of the polarizing films are changed.\nAnother approach to light control involves the use of polymer films or doping glass with metal ions to absorb or reject certain wavelength range of the electromagnetic radiation. Light transmission through windows using such technologies is fixed once the window structures are constructed.\nIn recent times, there has been great interest in using variable light transmission glass or glazing to achieve electromagnetic radiation control. A summary of variable tinting window developments can be found in, for example, “Chromogenic Switchable Glazing, Towards Development of the Smart Window” by Carl Lempert, published in the June 1995 Proceedings of the Window Innovation Conference, Toronto, Canada; and “Optical Switching Technology for Glazings” by Carl Lempert, published in Thin Solid Films, Volume 236, 1993, pages 6-13.\nSeveral different types of chromogenic switchable glazing structures have been discovered using suspended particle devices, electrochromic effects, and certain types of liquid crystal. In general, the structures absorb the incident light or diffuse the incident light.\nA glazing structure disclosed in U.S. Pat. No. 5,940,150 proposes an adjustable trans-reflective window that is primarily based on two layers of cholesteric liquid crystals and their ability to reject or to transmit one circular polarized light depending on the electric stimulus applied while capable of transmitting the other. However, a specific cholesteric liquid crystal typically works in only a narrow bandwidth of light and a stack of liquid crystal materials are needed to broaden the working wavelength range of the device. Thus, significant material as well as device development effort is needed in order that the glazing structure be practical. The patent also teaches a glazing structure sandwiching a controllable retardation film between two linear reflective polarizers. Such a glazing structure will reflect theoretically 50% to 100%, and transmit 50% to 0%, of naturally polarized light depending on the external stimulus applied. The theoretical contrast ratio between the maximum reflection state and the minimum reflection state of the glazing structure employing linear reflective polarizers disclosed in that patent is therefore only 2:1, making it less desirable as an electrically controllable mirror.\nThe foregoing examples of the related art and limitations related therewith are intended to be illustrative and not exclusive. Other limitations of the related art will become apparent to those skilled in the art upon a reading of the specification and a study of the drawings."}
{"text": "This invention relates to computer graphics and, more specifically to a parametric geometric element definition and generation system.\nMany applications such as on-line advertisements or promotional graphics require rich image data. As well, better computer graphics capabilities are required by interactive Internet based 3D programming and gaming as well as by real time Internet based simulations. These advanced capabilities are required by advertising agencies, people who create content for the Internet, game developers and entertainment companies. Existing systems are limited because of the large file size associated with graphic, images or image manipulation, the long time required to download content, the lack of rich animatable geometries, the long amount of time required to download level of detail (LOD) data, and the extensive memory and processing time required to generate adaptive level of detail information.\nIn a computer graphics system, graphic images may be created and/or manipulated by the computer. A graphic image may be displayed on a video display or output to a printer, for example. A graphic image is comprised of a plurality of geometric shapes, or polygons, and textures. As the level of detail increases, the number of polygons and/or textures in the graphic image and, consequently, the size of the image increases. The size of the graphic image affects the time that it takes for the image\"\"s display to be updated on the screen as well as the length of time it takes to send the image from one computer to another. It is desirable to be able to minimize the time needed to transmit an image and/or update its display.\nA display screen is comprised of a number of horizontal lines (referred to as scan lines) each of which is comprised of a number of dots (or pixels). To update an image on the screen, the area of the screen that contains the image is redisplayed by refreshing each pixel. This is typically performed by starting at the top left corner of the display and ending with the bottom right corner of the display. A complete scan of the display area is referred to as a frame. The number of frames that are displayed per second is referred to as the frame rate. The higher frame rate is typically desirable.\nThe frame rate is especially important where the graphic image is changed slightly and redisplayed to give the appearance of movement or animation. A low frame rate causes the animation to appear jerky. By lowering the level of detail (e.g., the number of polygons that must be drawn), the frame rate may be increased. However, a lower level of detail may cause the image to appear blocky and unrealistic.\nUsing existing approaches, each version of the image and/or different levels of detail of each version are created as different images (i.e., different xe2x80x9csnapshotsxe2x80x9d of the image). Before these images may be transmitted to another computer, a data representation of the images must be generated.\nOne mechanism for generating a data representation of an image is the Virtual Reality Modeling Language (VRML). VRML is a text-based definitional language comprising a statement syntax used to define a graphic image. A VRML file is created to define an image containing VRML statements that define nodes (e.g., sphere, cylinder, cone, cube, translation and transform) and fields (e.g., radius, height, width, and depth) that distinguish one node from another. The following is an example of VRML statements that define a red sphere:\nTransform {\ntranslation 301\nchildren [\nShape {\ngeometry Sphere {radius 2.3}\nappearance Appearance {material Material {\ndiffuseColor 100 #Red\n}}\n}\n]\n}\nThe VRML statements identify five nodes (i.e., xe2x80x9cTransformxe2x80x9d, xe2x80x9cShapexe2x80x9d, xe2x80x9cSpherexe2x80x9d, xe2x80x9cAppearancexe2x80x9d and xe2x80x9cMaterialxe2x80x9d) and fields associated with the nodes. The xe2x80x9cSpherexe2x80x9d node identifies a geometric shape (i.e., sphere) and the xe2x80x9cradiusxe2x80x9d field indicates that the sphere is to have a radius of 2.3. The xe2x80x9cAppearancexe2x80x9d and xe2x80x9cMaterialxe2x80x9d nodes affect the way the sphere is drawn. The material node with a xe2x80x9c1 0 0xe2x80x9d value for its xe2x80x9cdiffuseColorxe2x80x9d field indicates that the sphere is red. The xe2x80x9cShapexe2x80x9d node associates the xe2x80x9cAppearancexe2x80x9d node with the xe2x80x9cSpherexe2x80x9d node. The xe2x80x9cTranslationxe2x80x9d node defines a movement of the sphere to a new location as indicated by the xe2x80x9ctranslationxe2x80x9d field and its value of xe2x80x9c3 0 1xe2x80x9d. A VRML representation is read and the node and field information is used to generate an image.\nA VRML representation of an image represents a static state, or snapshot, of the image. Changes to the level of detail within the image, for example, are done by generating a new VRML representation of the image. Where the sphere is to be displayed on another computer, each of these representations must be sent to the computer before they can be used to display the sphere on the other computer.\nThe VRML representations of the sphere may be sent via the Internet, for example. Most users connect to the Internet using a modem with a maximum transmission rate that is frequently not achieved due to the quality of the phone line(s) used to transmit the data. The transmission rates using a modem connection are not sufficient to provide continuous streams of image data for processing by a computer. It would be beneficial to be able to the processing capabilities available with existing computer systems and minimize the amount of image data that must be transmitted between computer systems.\nU.S. Pat. No. 5,581,672 Dec. 3, 1996 Letcher, Jr. and U.S. Pat. No. 5,856,828, Jan. 5, 1999 Letchter, Jr. describe a system of relational entities for object-oriented computer-aided geometric design. In the patents, geometric objects are defined and constructed in a system which retains and utilizes relationsip and dependencies between objects in a directed graph data structure.\nA parametric geometric element definition and generation system and method is described herein. Embodiments of the invention comprise a model having a plurality of modeling operations and parameters or properties used in defining and generating geometric elements. In one or more embodiments of the invention, the model is a directed acyclic graph of three-dimensional modeling operations that combine to store information about a process used to construct an image. In one or more embodiments of the invention an object player monitors a model, communicates initialization and modifications of geometric data to the rendering system, and provides geometric data to the rendering system in the appropriate data formats.\nIn one or more embodiments of the invention, the model and modeling operations are object-oriented objects having properties and methods. One embodiment of the invention comprises geometry, core, operation and player object classes.\nGeometry object classes comprise the geometric data (e.g. coordinate, color, normal and texture coordinate data). The geometry object allows data to be shared by describing faces as lists of indices into the data pools. Core classes provide the fundamental communication and generation mechanism used by the model and comprise exceptions, change events, and operation superclasses. The operation classes comprise classes that may generate a geometric object or combine geometric objects to create new geometric objects. The object player provides an interface between the model and the underlying rendering system and may initiate a generation request. Geometric object data is returned to the object player in response to a generation request. This geometric object data is then forwarded to the underlying rendering system for rendering the image.\nA notification mechanism is used where objects may send and/or listen for change events. A modification to one or more properties of the model and/or modeling operations may trigger the transmission of a change event. The change event indicates to interested parties (e.g., other modeling operations, the model and for the object player) that an update to their geometric output, or other action, may be needed. A tracking mechanism provides the ability to determine what properties have been changed and assists in determining a response to a generate request. Each modeling operation has the ability to respond to a generate request. In an embodiment of the invention, the generation process is a depth-first traversal of the model graph structure. Generation may be limited to those operations of the model whose properties have been modified, and those operations that depend on the modified operation whose properties have been modified.\nDifferent xe2x80x9csnapshotsxe2x80x9d of an image may be generated at runtime by modifying one or more parameters of operations within a model of the image and then asking the graph to generate. Each snapshot may provide a different level of detail or animation, for example. In one embodiment of the invention, a predictive mechanism is used to determine an optimal level of detail.\nIn one or more embodiments of the invention, a prediction scheme is used to determine the effect on rendering time of image changes. Rendering times are collected and any differences in rendering times between rendering passes are allocated to the changed geometric objects. It is therefore possible to collect statistics on the effect on rendering time of making various changes to the image. It is possible, for example, to predict the rendering times for different levels of detail associated with a geometric object. Predictions may be used in selecting the appropriate level of detail to achieve a desired frame rate.\nIn accordance with a further aspect, the invention is directed to a method of computer aided 3D modelling of geometric objects undergoing changes wherein a directed acyclic graph is executed to generate a 3D model, the directed acyclic graph including a plurality of nodes which are connected with other nodes with directed edges between two nodes, each node representing a modelling operation of the geometric objects and each directed edge indicating the direction of flow of data from one node to another. The method comprises steps of a first node sending a change event indication specifying changes to bound properties of the node to a series of nodes in succession along associated directed edges and each node in succession determining if the change event indication received from a preceding node affects its bound properties and further determining based on its own rules that its own change event indication should be sent to a succeeding node. The method includes further steps of sending the change event indication to the succeeding node, and last node in succession instructing a rendering system to render a 3D model in accordance with the received change event indication.\nIn accordance with a yet another aspect, the invention is directed to a method of computer aided 3D modelling of geometric objects undergoing changes wherein a directed acyclic graph is executed to instruct a rendering system to generate a 3D model, the directed acyclic graph including a plurality of nodes which are connected with other nodes with directed edges between two nodes, each node representing a modelling operation of the geometric objects and each directed edge indicating the direction of flow of data from one node to another. The method comprises steps of (a) executing continually the directed acyclic graph to instruct the rendering system to generate a series of frames of a 3D model of geometric objects which are undergoing changes and (b) measuring a time required for the rendering system to generate a predetermined numbers of preceding frames with respect to a target time. The method further includes steps of (c) modifying bound properties of one or more nodes and (d) repeating iteratively steps (b) and (c) so that the time required for the rendering system to generate the most recent frame approach the target time.\nAccording to another aspect, the invention is directed to a computer aided 3D modelling apparatus for modelling 3D models of geometric objects undergoing change which comprises an instruction storage module containing a directed acyclic graph to be executed to generate a 3D model, the directed acyclic graph including a plurality of nodes which are connected with other nodes with directed edges between two nodes, each node representing a modelling operation of the geometric objects and each directed edge indicating the direction of flow of data from one node to another. The apparatus further includes an input port for inputting changes in bound properties, defined below, of one or more nodes, an instruction execution module for executing the directed acyclic graph to generate an output, an object player module for monitoring the output and regenerating the directed acyclic graph in response to the output if needed, and the rendering system for rendering the 3D model of geometric objects in response to the instructions."}
{"text": "Content displayed on a web page can be generated by one or more content item servers in response to content item requests that are generated during the rendering of the web page by a client device. Often a publisher of a web page may define multiple content item environments, e.g., advertisement slots, for content items, e.g., advertisements, to be displayed within a web page. Typically, each content item environment has an associated content item request e.g. an advertisement request, that is transmitted to a content item server, e.g., an advertisement server. When the web page is rendered by a client device that has requested the publishier's web page, a content item will be displayed in each content item environment.\nA web page having multiple content item environments can, however, generate excessive traffic to the content item server. Additionally, because web page instructions are often interpreted, the content item requests are generated as each content item request instruction is processed. Accordingly, a first content item request for a web page is transmitted and served without knowledge of any subsequent content item requests that may be encountered for that page. Thus, each content item served by the content item server for a web page if the rendering of the web page generates multiple content item requests."}
{"text": "Field of the Invention\nThe present invention relates to an imaging lens which forms an image of an object on a solid-state image sensor such as a CCD sensor or a C-MOS sensor used in a compact image pickup device, and more particularly to an imaging lens built in an image pickup device mounted in an increasingly compact and low-profile smartphone or mobile phone, PDA (Personal Digital Assistant), game console or information terminal such as a PC, or a home appliance with a camera function.\nDescription of the Related Art\nIn recent years, there has been a general tendency that mobile terminals such as smartphones and home appliances have a camera function. Today, whereas high-end models with a high-resolution camera function comparable to a digital still camera have been introduced into the market, the need for inexpensive popular models with a certain level of camera performance is still high. An imaging lens composed of three elements (constituent lenses) may be suitable as an imaging lens to be built in an inexpensive popular model because it can deliver a certain level of imaging performance and can be applied to a low-profile device and supplied at low cost. However, with the recent trend toward smaller higher-pixel image sensors, the pixel size is becoming smaller and the pixel density is becoming higher. Nowadays, image sensors with a pixel pitch of less than 1.2 microns have been proposed. An imaging lens for use in such an image sensor is expected to not only feature the smallness of aberrations but also provide a high-brightness optical system with an F-value smaller than 2.8 (often seen in the past). As for mobile terminals in particular, the imaging lens is expected to be low-profile enough to be applicable to a low-profile product. Furthermore, the imaging lens is anticipated to be able to capture an image of an object over a wide field of view so as to flexibly cope with the various camera functions of the product.\nConventionally, as an example of an imaging lens composed of three constituent lenses, the imaging lens described in JP-A-2010-113306 (Patent Document 1) includes, in order from an object side, an aperture stop, a first lens with positive refractive power, a second lens with negative refractive power, and a third lens. The third lens has, on the both sides, aspheric surfaces which are convex-curved toward the object near an optical axis and concave-curved toward the object in the vicinity of the lens periphery so that the refractive power changes according to the distance from the optical axis. The thickness of the first lens on the optical axis, the air gap on the optical axis between the first and second lenses, the focal length of the third lens, and the relation in curvature radius between the object-side surface and the image-side surface of the first lens are determined so as to achieve a wide field of view.\nThe imaging lens described in JP-A-2007-127953 (Patent Document 2) is an optical system which corrects chromatic aberrations using a diffractive optical surface. The imaging lens described in Patent Document 2 includes at least three constituent lenses and at least one of the lens surfaces of these lenses has a diffractive surface. At least one of the lens surfaces of the lens located nearest to the image plane has negative optical power in its center and the optical power changes to positive power as the distance to the lens periphery decreases, making up a compact optical system with high telecentricity.\nAccording to Patent Document 1, the imaging lens can capture an image of an object over a wide field of view from 76 to 78 degrees and corrects aberrations relatively properly. However, since the F-value is 2.8, there is difficulty in applying the imaging lens to a compact high-density image sensor as mentioned above. The imaging lens described in Patent Document 1 has a total track length of 4.6 to 5.5 mm and the ratio of the total track length to the diagonal length of the effective imaging plane of the image sensor (divide total track length by diagonal length of the effective imaging plane of the image sensor; hereinafter referred to as the ratio to diagonal) is about 1.0; however, in order to make the imaging lens more low-profile and offer brightness with an F-value of 2.8 or less, the problem related to correction of aberrations has to be addressed.\nAccording to Patent Document 2, the imaging lens can capture an image of an object over a maximum field of view of 72 degrees and corrects aberrations relatively properly. However, the positive, positive and negative refractive power lenses are arranged in order from the object side and chromatic aberrations are corrected not by the combination of the first and second lenses but by the diffractive optical surface. When the dependence on the diffractive optical surface for correction of chromatic aberrations is high, the number of orbicular zones formed on the lens surface tends to increase, which means that flare is more likely to occur. If two diffractive optical surfaces are formed to address this problem, high precision is required in the lens forming and assembling processes to prevent misalignment between the diffractive optical surfaces, which implies a higher degree of manufacturing difficulty. In connection with low-profileness, the third lens, located nearest to the image plane, has negative refractive power and the back focus is relatively long, making it difficult to achieve low-profileness. In the imaging lens composed of three constituent lenses as described in Patent Document 2, the total track length is about 6.8 mm and the ratio to diagonal is more than 1.1. If the back focus is decreased to achieve low-profileness, the angle of incidence on the third lens must be increased and as a consequence, probably the effective diameter of the second lens must be larger. When the lens diameter is larger, spherical aberrations and coma aberrations increase. This will make it difficult to make the imaging lens more low-profile while ensuring high performance. In addition, since the F-value of the imaging lens disclosed in Patent Document 2 is 3.3, it is difficult to apply the imaging lens to the latest high-density image sensors."}
{"text": "Currently, with the development of computer technology and electronic technology, electronic devices such as smart mobile phones, mobile Internet devices (MIDs), personal digital assistants (PDAs), and tablet computers have become indispensable tools for communication and entertainment in people's daily work and life. Users often use electronic devices for daily reading. When finding a content of interest during reading, users often want to make duplicates for sharing the content of interest. However, when the duplicated content includes both text information and pictures, the pictures in the duplicated content are often filtered out, causing the duplicated content inconsistent with the originally displayed content and affecting the user experience."}
{"text": "The present invention relates to an optical multi-wavelength selective device, more precisely a multi-wavelength selective switch for rerouting and switching of wavelengths in an optical network, in which each wavelength may be separately rerouted. The invention also comprises a method for rerouting and switching or add/drop multiplexing of wavelengths in an optical network.\nIn order to further increase the capacity on existing optical networks a number of methods are known. One means is to use so called wavelength division multiplexing (WDM) technique in order to improve the degree of utilization of the available bandwidth. In order to increase the network flexibility, devices able to reroute traffic in the optical network are necessary. Such devices for rerouting of the traffic are also suitable in order to utilize the network as efficiently as possible, and in case of network interrupts.\nIn xe2x80x9cWavelength Division Multiplexer with Photoinduced Bragg Gratings Fabricated in a Planar Lightwave Circuit Type Asymmetric Mach Zehnder Interferometer on Sixe2x80x9d Y. Hibino et.al., IEEE Photonics Technology Letters, Vol. 8, No. 1, January 1996 an optical element is described in which Bragg grating and a phase control element is used in a Mach Zehnder interferometer. The considered applications are wavelength division multiplexing and wavelength division de-multiplexing.\nThe above described device can however not be used as a multi-wavelength selective switch. If the described device is to be used for add/drop of several channels, as many devices are needed as the number of handled add/drop pairs. This type of device is comparatively hard to reconfigure, i.e. it is a non-adaptable device.\nIn K. W. Cheung, D. A. Smith, J. K. Baran, J. J. Johnson: xe2x80x9cWavelength-selective circuit and packet switching using acousto-optic tunable filtersxe2x80x9d Globecom \"\"90, pp 1541-1547, 1990, a manner is described of solving multi-wavelength selective rerouting of wavelength channels separate and independently of each other based on opto-aucoustic switches. This type of rerouting typically uses a number of micro-seconds, which may be a problem in high speed transmissions.\nIn R. C. Alferness, L. L. Buhl, U. Koren, B. I. Miller, M. Young, T. L. Koch.xe2x80x9d Broadly tunable InGaAsP/InP buried rib waveguide vertical coupler filterxe2x80x9d, Techn. Dig. Integrated Photonics Research, post-deadline paper, pp. 33-36, 1991 a tunable wavelength selective grating-assisted directional coupler comprising two waveguides having different effective refractive indexes and a grating described. This directional coupler can couple a wavelength from the first waveguide to the second waveguide.\nThe problem present in the known technique as regards multi-wavelength selective switches and re-directional devices is this demand very big and complicated structures or more components, which results in big power losses and high costs, or in that they can not be used in connection with high speed transmissions.\nIn order to increase the capacity in an optical transmission system several know methods may be used. In wavelength division multiplexing the transmission channels are multiplexed and de-multiplexed on different carrier wavelengths to respectively from an information flow. Said multiplexing and de-multiplexing demands optical multi-wavelength selective devices. It may also be desirable to change the transmission route through the optical net for certain wavelength channels. For this purpose components are needed having multi-wavelength selective performance characteristics such as a multi-wavelength selective switch.\nA known problem related to multi-wavelength selective switches may be that these contribute to big power losses.\nAnother problem is that all by us known multi-wavelength selective switches known to us exhibited a relatively complicated structure or a great number of components.\nAnother problem is that known multi-wavelength selective switches are relatively expensive to manufacture depending on said complicated structures or the number of components comprised.\nYet a problem is that certain known multi-wavelength selective switches exhibit relatively slow switching speeds, typically a number of microseconds.\nThe present invention approaches the said problem through a multi-wavelength selective switch in the form of a directional coupler structure comprising two waveguides arranged having different effective refractive index. The waveguides are at one section arranged in the proximity of each other, such that an optical field in each of the waveguides may overlap the other wherein at least one grating is arranged for co-directional coupling.\nIn a method according to the invention for switching and re-direction of optical wavelengths in an optical network one branch in a directional coupler structure comprising optical wavelengths is excited. The optical wavelengths are thereafter transmitted through said branch to the directional coupler structure. The optical wavelengths are thereafter transmitted through a section where the waveguides are arranged in the proximity of each other whereby none or at least one wavelength is influenced by at least one in this section arranged grating for co-directional coupling and is coupled to the other waveguide.\nIn a second embodiment according to the invention at least one wavelength selective amplifier is comprised in at least one branch of the directional coupler structure, e.g. in the form of a distributed-Bragg-reflector-amplifier.\nThe object of the present invention is to obtain a multi-wavelength selective switch, the structure of which is so simple and compact as possible, such that the manufacturing thereby may be reduced compared to the available technique, and in which the switching speed still is relatively high.\nAn advantage of the present invention is that the power losses relatively seen may be kept low.\nA further advantage of the present invention is that its remaining performance characteristics such as crosstalk and the like may be improved compared to know technique.\nStill a further advantage of the present invention is that signal splitting only is performed when absolutely necessary.\nThe invention will now be described further using the preferred embodiments and with reference to the appended drawings."}
{"text": "This invention relates to subsea wellhead systems and more particularly, to methods and apparatus for supporting, holding down, and sealing casing hangers within a subsea wellhead.\nIncreased activity in offshore drilling and completion has caused an increase in working pressures such that it is anticipated that new wells will have a working pressure of as high as 15,000 psi. To cope with the unique problems associated with underwater drilling and completion at such increased working pressures, new subsea wellhead systems are required. Wells having a working pressure of up to 15,000 psi are presently being drilled off the coast of Canada and in the North Sea in depths of over 300 feet. These drilling operations generally include a floating vessel having a heave compensator for a riser and drill pipe extending to the blowout preventer and wellhead located at the mud line. The blowout preventer stack is generally mounted on 20 inch pipe with the riser extending to the surface. A quick disconnect is often located on top of the blowout preventer stack. An articulation joint is used to allow for vessel movement. Two major problems arise in 15,000 psi working pressure subsea wellhead systems operating in this environment, namely, a support shoulder in the wellhead housing which will support the casing and pressure load, and a sealing means between the casing hangers and wellhead which will withstand and contain the working pressure.\nIn the past, prior art wellhead designs permitted adequate landing support for successive casing hangers. However, with the increase in pressure rating and the landing and supporting of multiple casing strings and tubing strings within the wellhead, a small support shoulder will not support the load. Although an obvious answer to the problem would be to merely use a support shoulder large enough to support the casing and pressure load, large support shoulders projecting into the flow bore in the wellhead housing restrict access to the casing below the wellhead housing for drilling. In the early days of offshore drilling, 16-3/4 inch bore subsea wellhead systems required underreaming. At that time, most floating drilling rigs were outfitted with a 16-3/4 inch blowout preventer system to eliminate the two stack (20 inch and 13-5/8 inch) and the two riser system required up until that time. As wellhead systems moved from 5,000 psi to 10,000 psi working pressure, the 18-3/4 inch, 10,000 psi support shoulder was developed to carry casing and pressure loads and to provide full access into the casing below the wellhead housing.\nThe second major problem is the sealing means. The sealing means must be capable of withstanding and containing 15,000 psi working pressures. Available energy sources for energizing the sealing means include weight, hydraulic pressure, and torque. Each sealing means requires different amounts of energy to position and energize. Weight is the least desirable because the handling of drill collars providing the weight is difficult and time consuming on the rig floor. If hydraulic pressure is applied through the drill pipe, there is a need for wireline equipment to run and recover darts from the hydraulic-to-actuated seal energization system. If darts are not used, the handling of \"wet string\" of drill pipe is very messy and unpopular with drilling crews. If the seal energization means uses the single trip casing hanger technique, the cementing fluid can cause problems in the hydraulic system used to energize the seal. Maintenance is also a problem. Although torque is the most desirable method to energize a seal, there are limitations on the amount of torque which can be transmitted from the surface due to friction losses to riser pipe, the blowout preventer stack, off location, various threads, and the drill pipe itself.\nThe subsea wellhead system of the present invention overcomes the deficiencies of the prior art and includes many other advantageous features. The system is simple, has less than 50 parts and is suitable for H.sub.2 S service. The system has single trip capability but can still use multiple trip methods. All hangers are interchangeable with respect to the outer profile so that they can be run in lower positions. The seal elements are interchangeable and are fully energized to a pressure in excess of the anticipated wellbore pressure. Back-up seals are available. The seals are not pressure de-energized. The hangers can be run without lock downs and the seal elements will seal even if the hanger lands high.\nThe housing support seat supports in excess of 6,000,000 lbs. (working pressure plus casing weight or test pressure) without exceeding 150% of material yield in compression. The wellhead will pass a 17-1/2 inch diameter bit. The present invention does not attempt to land on two types of seats at once or on two seats at once. Further, the housing support seat is not sensitive to collecting trash during drilling or to collecting trash during the running of a 13-3/8 inch casing. Further, the housing support seat does not require a separate trip nor does it drag snap rings down the bore.\nThe hanger hold down will hold down 2,000,000 lbs. The hanger hold down is positively mechanically retracted when retrieving the casing hanger body and is compatible with single trip operations. The hanger hold down is released for retrieval of the casing hanger when the seal element is retrieved. The hanger hold down is compatible with multiple trip operations and permits the running of the hanger with or without the hold down. The sealing means will work even if the hold down is not used. The hanger hold down is reusable and has a minimum number of tolerances to stack up between hold down grooves.\nThe sealing means of the present invention will reliably seal an annular area of approximately 18-1/2 inch outside diameter by 17 inch inside diameter and provide a rubber pressure in excess of 15,000 psi (20,000 psi nominally) when the sealing means is energized and the sealing means sees a pressure from above or below of 15,000 psi. The pressure in excess of 15,000 psi is retained in the sealing means after the running tool is removed. The sealing means is additionally self-energized to hold full pressure where full loading force was not applied or where full loading force was not retained. The sealing means will not be pressure de-energized. The sealing means provides a relatively long seal area to bridge housing defects and/or trash. Further, the sealing means provides primary metal-to-metal seals and used the metal-to-metal seals as backups to prevent high pressure extrusion of secondary elastomeric seals. The sealing means of the present invention positively retracts the metal-to-metal seals from the walls prior to retrieving the sealing means. The elastomeric seals of the sealing means are allowed to relax during retrieval of the packoff assembly and is completely retrievable. The present sealing means provides a substantial metallic link between the top and the bottom of the packing seal area to insure that the lower ring is retrievable. The design allows for single trip operations. There are no intermittent metal parts in the seal area to give irregular rubber pressures. The sealing means provides a minimum number of seal areas in parallel to minimize leak paths. The sealing means is positively attached to the packing element so that it cannot be washed off by flow during the running operations. The design also allows for multiple trip operations and is interchangeable for all casing hangers within a nominal size.\nThe means to load the sealing means reliably provides a force to energize the sealing means to a nominal 20,000 psi. It allows full circulation if used in a single trip. However, the loading means is compatible with either a single trip operation or multiple trip operation. Further, it is interchangeable for all casing hangers within the wellhead system. The loading means will cause the sealing means to seal even if the casing hanger is set high. Further, it does not release any significant amount of the full pressure load after actuation. The loading means does not require a remote engagement of hold down threads. Further, it has no shear pins. The loading means is reusable and does not have to remotely engage hold down threads on packing nut replacement.\nThe casing hanger running tool includes a connection between the running tool and casing hanger which will support in excess of 700,000 lbs. of pipe load. The running tool is able to generate an axial force in excess of 900,000 lbs. to energize the sealing means. Further, the running tool is able to tie back into the casing hanger without a left hand torque. The running tool can be run on either casing or drill pipe.\nOther objects and advantages of the invention will appear from the following description."}
{"text": "Chromatography and solid-phase extraction (“SPE”) are commonly-used separation techniques employed in a variety of analytical chemistry and biochemistry environments. Chromatography and SPE are often used for separation, extraction, and analysis of various constituents, or fractions, of a sample of interest. Chromatography and SPE may also be used for the preparation, purification, concentration, and clean-up of samples.\nChromatography and SPE relate to any of a variety of techniques used to separate complex mixtures based on differential affinities of components of a sample carried by a mobile phase with which the sample flows, and a stationary phase through which the sample passes. Typically, chromatography and SPE involve the use of a stationary phase that includes an adsorbent packed into a cartridge, column, or disposed as a thin layer on a plate. Thin-layer chromatography (“TLC”) employs a stationary phase that is spread in a thin layer on a carrier or substrate plate. A commonly-used stationary phase includes a silica-gel-based sorbent material.\nMobile phases are often solvent-based liquids, although gas chromatography typically employs a gaseous mobile phase. Liquid mobile phases may vary significantly in their compositions depending on various characteristics of the sample being analyzed and on the various components sought to be extracted and/or analyzed in the sample. For example, liquid mobile phases may vary significantly in pH and solvent properties. Additionally, liquid mobile phases may vary in their compositions depending on the characteristics of the stationary phase that is being employed. Often, several different mobile phases are employed during a given chromatography or SPE procedure. For example, gradient elution may be performed in which the mobile phase composition is varied with time.\nA typical TLC plate is prepared by mixing an adsorbent (which acts as the stationary phase) with a small amount of an inert binder and water. The mixture may be spread as a relatively viscous slurry onto a carrier sheet. The resulting plate can then be dried and activated in an oven. The resulting stationary phase is bound in place to the carrier sheet or other substrate by the binder. The presence of the binder can lead to secondary interactions with the mobile phase, as well as a decrease in separation efficiency."}
{"text": "A. Technical Field\nThe present invention relates to a process for producing a hydroxyalkyl (meth)acrylate by performing a reaction between (meth)acrylic acid and an alkylene oxide.\nB. Background Art\nIn the case of producing a hydroxyalkyl (meth)acrylate by performing a reaction between (meth)acrylic acid and an alkylene oxide, the formation of a diester (alkyleneglycol di(meth)acrylate) as a by-product has so far been a problem. The diester has two unsaturated bonds within one molecule, and promotes the polymerization of an object, the hydroxyalkyl (meth)acrylate, and causes the trouble such as clogging of an apparatus. Furthermore, in the case of producing a (co)polymer from a raw material of the hydroxyalkyl (meth)acrylate containing the diester, turbidity arises in the resulting polymer, or the hydroxyalkyl (meth)acrylate containing the diester causes the occurrence of gelation at the time of polymerization. Furthermore, the vapor pressure of the diester is approximate to that of the hydroxyalkyl (meth)acrylate. Therefore, once the diester is formed, the separation thereof is extremely difficult. In addition, when an object, the hydroxyalkyl (meth)acrylate is distilled and purified, the diester is also formed by a disproportionation reaction thereof.\nIt is known that the diester is apt to form when a reaction liquid of (meth)acrylic acid and the alkylene oxide, or the resulting hydroxyalkyl (meth)acrylate is elevated to a high temperature (e.g. JP-A-133227/1976). Therefore, if a reaction temperature or a distillation temperature is low, the formation of the diester is controlled to some degree. However, when the reaction temperature is low, the period of reaction step time lengthens very much, and so that there is not only a possibility that the productivity falls, but also a possibility that part of a product polymerizes because of the long period of reaction step time or that other by-products increase. Accordingly, falling the reaction temperature or the distillation temperature is not a practical countermeasure.\nA method of controlling the formation of the diester at the time of distillation by deactivating a specific catalyst used in the reaction is disclosed (e.g. JP-A-194321/1993). However, there is a disadvantage that the available catalyst is limited, or that the catalyst is not applicable to the formation of the diester at the time of the reaction, and the control effect thereof was insufficient.\nA. Object of the Invention\nAn object of the present invention is to provide a process for producing a hydroxyalkyl (meth)acrylate which enables to control the formation of a diester sufficiently at the time of reaction and/or at the time of distillation.\nB. Disclosure of the Invention\nThe present inventors diligently studied to achieve the above object. As a result, they have found that the above object could be solved completely by allowing a liquid with a specific relative dielectric constant to coexist at the time of a reaction between (meth)acrylic acid and an alkylene oxide and/or at the time of distillation of a hydroxyalkyl (meth)acrylate.\nThat is to say, a process for producing a hydroxyalkyl (meth)acrylate, which comprises the step of carrying out a reaction between (meth)acrylic acid and an alkylene oxide in the existence of a catalyst; with the process being characterized in that when the reaction between (meth)acrylic acid and the alkylene oxide is carried out, there coexists a liquid having a relative dielectric constant of not less than 20 at 25xc2x0 C.\nAn another process for producing a hydroxyalkyl (meth)acrylate, which comprises the step of carrying out a reaction between (meth)acrylic acid and an alkylene oxide in the existence of a catalyst; with the process being characterized by further comprising the step of distilling the resultant hydroxyalkyl (meth)acrylate in the coexistence of a liquid having a relative dielectric constant of not less than 20 at 25xc2x0 C.\nThese and other objects and the advantages of the present invention will be more fully apparent from the following detailed disclosure.\nFirst, an outline of a process for producing a hydroxyalkyl (meth)acrylate which the characteristic production process of the present invention can favorably apply to is explained as follows.\nFirst, a reaction between (meth)acrylic acid and an alkylene oxide is carried out in the existence of a catalyst. The reaction rate of this reaction does not reach 100% in many cases, and it is general that the unreacted (meth)acrylic acid or alkylene oxide or such remains in a reaction liquid at the time of the end of the reaction. Then, the above reaction liquid is introduced into the step of removing the unreacted raw materials and such from the reaction liquid. As the subsequent final step, purification such as distillation is performed to obtain an object, the hydroxyalkyl (meth)acrylate.\nHereinafter, the features of the production process of the present invention are explained below.\nThe alkylene oxide which can be used in the present invention are favorably an alkylene oxide having 2 to 6 carbon atoms, more favorably an alkylene oxide having 2 to 4 carbon atoms. Examples thereof include ethylene oxide, propylene oxide, and butylene oxide, more favorably ethylene oxide and propylene oxide. Incidentally, (meth)acrylic acid used in the present invention means acrylic acid or methacrylic acid.\nWhen the present invention is carried out, the amount of raw materials as charged for the above reaction between the (meth)acrylic acid and the alkylene oxide is favorably in the range of not less than 1 mol of the alkylene oxide per 1 mol of the (meth)acrylic acid, more favorably in the range of 1.0 to 5.0 mol, still more favorably in the range of 1.0 to 3.0 mol, especially favorably in the range of 1.0 to 2.0 mol. In the case where the amount of the alkylene oxide as charged is less than 1.0 mol, it is unfavorable because the reaction rate falls and by-products increase. When the amount of the alkylene oxide as charged is too much, especially more than 5 mol, it is unfavorable economically.\nIn the present invention, the reaction between the (meth)acrylic acid and the alkylene oxide in the existence of the catalyst can be performed by a method which is conventionally used for this kind of reaction.\nIn the case where the reaction is performed by a batch method, the reaction is performed by introducing a liquid alkylene oxide into the (meth)acrylic acid. The alkylene oxide may be introduced after the (meth)acrylic acid is dissolved in a solvent. In this case, the alkylene oxide may be added all at once, continuously, or intermittently. In the case of adding the alkylene oxide continuously or intermittently, as is often the case with this kind of reaction, the reaction can be completed by the so-called aging that is carried out by continuing the reaction after introducing the alkylene oxide. Furthermore, it is not always necessary to add the (meth)acrylic acid all at once in the initial stage, and the (meth)acrylic acid can be added after being divided into some portions.\nIn the case where the reaction is performed by a continuous method, the reaction is performed by adding the (meth)acrylic acid and the liquid alkylene oxide into a reactor such as a tubular or tank reactor continuously and extracting a reaction liquid out of the reactor continuously. In this case, a catalyst may be provided continuously together with raw materials and b e extracted continuously together with the reaction liquid. In the case of using the reactor such as the tubular reactor, a so-called fixed bed manner in which a solid catalyst as packed with the reactor is used may be utilized. In the case of using the reactor such as the tank reactor, the so-called fluidized bed manner in which a solid catalyst as fluidized together with the reaction liquid in the reactor is used may be used. Furthermore, in the cases of these continuous reactions, modes which allow part of the reaction liquid to circulate may be used.\nThe reaction temperature is usually favorably in the range of 40 to 130xc2x0 C., more favorably 50 to 100xc2x0 C. If the reaction temperature is below 40xc2x0 C., it takes more time for the reaction to proceed and it becomes apart from a practical level. On the other hand, if the reaction temperature is above 130xc2x0 C., it is unfavorable because by-products become large, or the polymerization of the (meth)acrylic acid of a raw material, the hydroxyalkyl (meth)acrylate of the product and such arise.\nFurthermore, this reaction may be performed in a solvent for the purpose of allowing the reaction to proceed mildly or such. As to the solvent, conventional ones such as toluene, xylene, heptane and octane are usable. The pressure in a reaction system during the reaction depends on the kinds or mixing ratios of used raw materials. However, the reaction is usually performed under applied pressure.\nIn the reaction, conventional polymerization inhibitors are usable. Examples thereof include hydroquinone, hydroquinone monomethyl ether, phenothiazine, copper dibutyldithiocarbamate, or such. The amount of the polymerization inhibitor as added is not especially limited, and is favorably in the range of 0.001 to 1 weight % of a carboxylic acid, more favorably in the range of 0.01 to 0.5 weight %.\nIn the production process of the present invention, a catalyst used in the reaction between the (meth)acrylic acid and the alkylene oxide is not especially limited, and a catalyst which is used conventionally in this kind of reaction can be used. Examples thereof include iron compounds such as ferric chloride and ferric acetate; chromium compounds such as chromium chloride, chromium acetylacetonate, chromium formate, chromium acrylate, chromium methacrylate, chromium dibutyldithiocarbamate, chromium nitrate, chromium sulfate; amines such as trialkylamines and ion-exchange resins having a quaternary ammonium group. It is favorable to use a basic catalyst such as an amine compound or a basic resin. This basic resin is a high-molecular compound with a basic functional group which is insoluble in a reaction liquid. Examples thereof include tertiary amine compounds, quaternary ammonium salts, cyclic amine compounds such as pyridine, high-molecular compounds with sulfides, favorably basic anion-exchange resins, especially favorably basic anion-exchange resins having an amino group as a basic functional group.\nIncidentally, the production process of the present invention enables to reveal the effect of controlling the formation of the diester sufficiently without using a chromium compound with valence number of 6 such as chromic anhydride which has so far exhibited the relatively high effect of controlling the formation of the diester as a catalyst.\nIn the production process of the present invention, distillation of the product of the hydroxyalkyl (meth)acrylate can be performed according to a method which is conventionally used for this kind of distillation, and for example a distillation apparatus having a portion of a vacant column is enumerated. Furthermore, a thin-film evaporation apparatus may be used in combination.\nIn the case of using the distillation apparatus with the portion of the vacant column, the so-called distillation column is favorable. The structure thereof can be simple distillation and examples thereof include rectifying columns such as a packed column, a bubble cap tower and a perforated-tray column. Especially, a distillation column comprising plural sieve trays or a distillation column comprising packed materials is favorable. Furthermore, a multi-plated perforated-tray column is desirable from a viewpoint of the maintenance of purity or the removal of a polymer. In the method of heating, a method which uses jacket type heaters or a reboiler method which allows bottoms to circulate and performs heating with shell-and-tube heat-exchangers is included. With respect to the pressure condition of distillation, it is favorable to perform distillation under reduced pressure. It is desirable that the operating pressure should be as low as possible from a viewpoint of preventing polymerization, favorably in the range of 0.66 to 40 hPa, more favorably in the range of 1.33 to 13 hPa, still more favorably in the range of 4 to 9.3 hPa. A liquid temperature in the distillation apparatus is usually favorably in the range of 70 to 130xc2x0 C., more favorably in the range of 80 to 110xc2x0 C.\nThe production process according to the present invention is characterized by carrying out the reaction between the (meth)acrylic acid and the alkylene oxide and/or distillation of the resultant hydroxyalkyl (meth)acrylate as explained in the above in the coexistence of a liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. This enables to exhibit the effect of controlling the formation of the diester sufficiently on the occasion of the reaction between the (meth)acrylic acid and the alkylene oxide and/or distillation of the hydroxyalkyl (meth)acrylate.\nThe liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. includes water, dimethyl sulfoxide, acetylacetone, acetaldehyde, acetone, acetonitrile, isobutyronitrile, ethanol, ethylene glycol, glycolonitrile, glycerin, o-chlorobenzonitrile, m-chlorobenzonitrile, succinonitrile, o-nitrotoluene, m-nitrotoluene, p-nitrotoluene, 1-nitropropane, 2-nitropropane, nitrobenzene, nitromethane, 1,4-butanediol, furfural, propiononitrile, benzonitrile, formamide, methanol, lactonitrile, and the like, preferably water, dimethyl sulfoxide, acetonitrile, furfural, benzonitrile and formamide, especially preferably at least one member selected from the group consisting of water, dimethyl sulfoxide and acetonitrile. In the case of a liquid having the relative dielectric constant of less than 20 at 25xc2x0 C., it is difficult to achieve the effect of the present invention, although the liquid is allowed to coexist at the time of the reaction between the (meth)acrylic acid and the alkylene oxide and/or distillation of the hydroxyalkyl (meth)acrylate.\nIn the case where a liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. is coexistent in a reaction liquid, the coexistence ratio of the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. per the reaction liquid is favorably in the range of 0.01 to 70 weight %, more favorably in the range of 0.1 to 50 weight %, especially favorably in the range of 1 to 40 weight %. When the coexistence ratio is less than 0.01 weight %, it is unfavorable because the formation amount of the diester increases, and when the coexistence ratio exceeds 70 weight %, it is unfavorable because the production efficiency deteriorates with the result that the production amount decreases or that it becomes necessary to enlarge production facilities.\nFurthermore, in the case where the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. is coexistent in a liquid which is fed to the distillation apparatus, the coexistence ratio of the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. per the liquid which is fed to the distillation apparatus is favorably in the range of 0.01 to 70 weight %, more favorably in the range of 0.1 to 50 weight %, especially favorably in the range of 1 to 40 weight %. When the coexistence ratio is less than 0.01 weight %, it is unfavorable because the formation amount of the diester increases to bring about the deterioration of product quality and the fall of product purity, and when the coexistence ratio exceeds 70 weight %, it is unfavorable because the construction cost of facilities involved in the enlargement of the distillation apparatus increases.\nIncidentally, a partial or full catalyst used for the reaction between the (meth)acrylic acid and the alkylene oxide may remain in the liquid which is fed to the distillation apparatus.\nIn the present invention, a mode that the reaction between the (meth)acrylic acid and the alkylene oxide is performed in the coexistence of the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C., or a mode that distillation of the hydroxyalkyl (meth)acrylate is performed in the coexistence of the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. is not especially limited, and examples thereof include (1) a mode that the reaction is performed after adding the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. on the occasion of adding reaction raw materials, (2) a mode that the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. is not coexistent in the reaction and distillation is performed after adding the concerned liquid in distillation, and (3) a mode that the reaction is performed after adding the liquid having the relative dielectric constant of not less than 20 at 25xc2x0 C. on the occasion of adding reaction raw materials and furthermore distillation is performed in the coexistence of the concerned liquid as it is.\n(Effects and Advantages of the Invention)\nA process for producing a hydroxyalkyl (meth)acrylate, according to the present invention, enables to control sufficiently the formation of a diester at the time of a reaction and/or at the time of distillation."}
{"text": "The following abbreviations and terms are herewith defined, at least some of which are referred to within the following description of the state of the art and the present invention.\n3GPP 3rd Generation Partnership Project\nACL Access Control List\nBFD Bidirectional Forwarding Detection\nCS Circuit Switch\nDA Destination Address\nICMP Internet Control Message Protocol\nIMS IP Multimedia Subsystem\nIM-MGW IP Multimedia Media Gateway Function\nIP Internet Protocol\nIPLx IP Level x\nIPv4 IP version 4\nIPv6 IP version 6\nISO International Organization for Standardization\nLSP Label Switch Path\nMGW Media Gateway\nMPLS Multi-Protocol Label Switching\nMSCS Mobile-Services Switching Center\nNHR Next Hop Router\nNUD Network Unreachability Detection\nOSI Open System Interconnect\nRFC Request for Comment\nSA Source Address\nTS Technical Specification\nTTL Time To Live\nNext Hop Router (NHR): The first router in the path from the host to the destination network.\nIMS Network—IPv4 and IPv6\nThe IP Multimedia Subsystem (IMS) originally specified the use of the new IP version 6 (IPv6) technology and later allowed the legacy IP version 4 (IPv4), but the intent is to move forward with IPv6 communications. Hence, the network providers are requiring IPv6 interfaces with their IMS networks. However, the access nodes from the previous non-IMS technologies, e.g.:                Call Server, e.g., Mobile-Services Switching Center (MSCS)        Media Gateway (MGW)may or may not support native IPv6, and often do not. Thus, the CS's MSCS is having capabilities shown as MGCF added so it can interface with the IMS network. Likewise, the CS's MGW is having capabilities shown as IM MGW added so it can support the IMS Mn and Mb interfaces. FIG. 1 (PRIOR ART) is a diagram that illustrates the logical interfaces CS 102a and 102b respectively between the IMS network's IMS-MGW 104 and MGCF 106 and the CS network's MGW 108 and MSCS 110—this diagram is part of TS 23.228 3GPP; TS Group Services and System Aspects; IP Multimedia Subsystem (IMS); Stage 2 (Release 11), 2011-12 (the contents of which are incorporated herein by reference). In this case, IMS components, such as the Serving-Call Session Control Functions (S-CSCF) 112 and the Media Resource Function Processor (MRFP) 114 respectively communicate with the MGCF 106 and the IMS-MGW 104 using the IPv4 and/or IPv6 (preferred) protocols. These IP capable host devices—S-CSCF 112/MGCF 106 and MRFP 114/IMS-MGW 104—will communicate with each other over an IP capable network infrastructure composed of IP routers which are not shown in FIG. 1 because they are associated with physical connections rather than logical connections.        \nThe IMS-MGW 104 and the MGCF 106 both require IPv4 and IPv6 capabilities and so that these devices can utilize the existing direct-connected IPv4 infrastructure they both implement an IPv4 to IPv6 transition technique known as tunneling. FIG. 2 (PRIOR ART) illustrates the IMS-MGW 104 and MGCF 106 which function as hosts interacting with their next hop routers 202 and 204 which are associated with the IMS subsystem IPv6 206 utilizing tunneling where IPv6 is tunneled over IPv4 so that the existing host IPv4 infrastructure does not change and the IPv6 is managed at the host's application end point. The next hop routers 202 and 204 are part of the aforementioned IP network infrastructure. The next hop routers 202 and 204 are capable of tunneling over IPv4 the native IPv6 traffic between the IMS components—S-CSCF 112/MGCF 106 and MRFP 114/IMS-MGW 104.\nThe Internet Protocol (IP) may utilize tunneling techniques to transport information for various reasons:                Security        IPv4 to IPv6 Transition        IPv6 to IPv4 TransitionPacket Encapsulation        \nFIG. 3 (PRIOR ART) illustrates how tunneling is realized by encapsulating a packet 300 with its original IP header 302 and data 304 within a new IP header 306 to form an encapsulated packet 308. The IP headers 302 and 306 may be any combination of IPv4 or IPv6. For instance, the original IP header 302 is an IPv6 IP header and the new IP header 306 is an IPv4 IP header, or the original IP header 302 is an IPv4 IP header and the new IP header 306 is an IPv6 IP header, or the original IP header 302 is an IPv4 IP header and the new IP header 306 is an IPv4 IP header, or the original IP header 302 is an IPv6 IP header and the new IP header 306 is an IPv6 IP header\nTunnel Endpoints\nFIG. 4 (PRIOR ART) illustrates a tunnel 400 (which includes the encapsulated packet 308) with two tunnel end points 402a and 402b each represented by an IP address associated with each IP layer. For example, an IPv6 over IPv4 tunnel 400 would have both an IPv6 address (IP Level 2) 404a and 404b and an IPv4 address (IP Level 1) 406a and 406b at each tunnel endpoint 402a and 402b. In particular, the tunnel endpoint 402a would have an IP address L2a 404a and an IP address L1a 406a while the tunnel endpoint 404b would have IP address L2b 404b and IP address L1b 406b. The IP addresses 406a and 406b are used to route between a host (e.g., MGCF 106) and a router (e.g., router 204) as they are used in the outside new IP header 306 And, the IP addresses 404a and 404b are used to route between the host and the far end device (see FIG. 6). The word “level” is used herein to distinguish the different stacked IP headers and is not to be confused with the term “layer” which is used in the International Organization for Standardization (ISO) Open System Interconnect (OSI) 7-layer Reference Model.\nTunnels vs. Links\nTunnels may exist between a host and routers (host-routers) or between the routers themselves (router-router). An algorithm for tunneling IPv6 over IPv4 between a host-and-routers and between routers-and-routers is described in RFC 4213 “Basic Transition Mechanisms for IPv6 Hosts and Routers” October 2005 (the contents of which are incorporated by reference herein). The IP protocol operates at the Network Layer, or Layer 3, of the ISO OSI Reference Model. This implies that the tunnel operation is also at Layer 3.\nAs shown in FIG. 5 (PRIOR ART), the host 500 when connecting to next hop routers 5021 and 5022 (only two shown) usually utilizes redundant paths 5061 and 5062 to the adjacent routers 5021 and 5022 to ensure connectivity even if there is a single point of failure (see also FIG. 2). The host 500 has an algorithm 510 which determines which link or links are utilized when communicating via its adjacent routers 5021 and 5022 to the far end device 504. This algorithm 510 will therefore determine the physical path(s) 5061 and 5062 used for communicating via the adjacent router(s) 5021 and 5022 to the far end device 504. Note: the physical paths 5061 and 5062 along with their associated data link protocol operate at the Data Link (Layer 2) and the Physical (Layer 1) of the OSI model and are therefore segregated from the tunneling which operates at Layer 3.\nThe host 500 to survive a single point of failure needs a minimum of two paths 5061 and 5062 connected to the two routers 5021 and 5022. Although more paths and routers may be utilized, the two-path algorithm is generally deployed. The adjacent routers 5021 and 5022 generally have an inter-router link 512 (the link inter-connecting each of the adjacent routers 5021 and 5022) to provide alternative path routing for both ingress traffic and egress traffic with respect to the host 500 and the network (far end device 504).\nAs shown in FIG. 6 (PRIOR ART), there is a host-router configuration in which the host 500 inter-connects to its adjacent or next hop router 5021 (for example) using one network protocol “A”, e.g., IPv4, while an application 602 on the host 500 needs to communicate with another application 606 on the far end device 504 using a different network protocol B, e.g., IPv6. Thus, to allow the host 500 to communicate with the far end device 504 there is established a host-router tunnel 6041 which uses network protocol “A” with network protocol “B” embedded therein.\nAs shown in FIG. 7 (PRIOR ART), if a host-router tunnel is utilized, then the host 500 will need a tunnel 6041, 6042 . . . 604n to each of its adjacent routers 5021, 5022 . . . 502n. Note: that a tunnel, and its traffic, between two end points may utilize different physical paths over time.\nTunneling and Routing\nHost-to-Router Direction\nIn the host 500 to router 5021 (for example) packet flow direction, the original IP header contains:                Destination IP Address—Far end device's address        Source IP Address—Host's address        \nThe host 500 will encapsulate the original IP packet with a new IP header. The new IP header contains:                Destination IP Address—Router's Tunnel Endpoint Address        Source IP Address—Host's Tunnel Endpoint Address        \nOnce the packet is removed from the tunnel by the router 5021, the external IP header is removed and the original IP header is then used to route the packet.\nRouter-to-Host Direction\nIn the router 5021 (for example) to host 500 packet flow direction, the original IP header contains:                Destination IP Address—Host address        Source IP Address—Originator of the IP packet (far end device 504)        \nThe router 5021 will encapsulate the original packet with the tunnel header. The new IP header contains:                Destination IP Address—Host's Tunnel Endpoint Address        Source IP Address—Router's Tunnel Endpoint AddressTunnel Integrity Check        \nTunnels 6041, 6042, 604n (for example) are effectively managed by software and are therefore subject to failure. The host 500 (for example) could check the tunnel's integrity by running heartbeat messages at the tunnel level. Examples of such a heartbeat mechanism 608 include:\nBidirectional Forwarding Detection (BFD):\n                RFC 5880—Bidirectional Forwarding Detection (BFD), Standard, June 2010.        RFC 5881—Bidirectional Forwarding Detection (BFD) for IPv4 and IPv6 (Single Hop), June 2010.        RFC 5883—Bidirectional Forwarding Detection (BFD) for Multihop Paths, June 2010.        RFC 5884—Bidirectional Forwarding Detection (BFD) for MPLS Label Switched Paths (LSPs), June 2010.        RFC 5882—Generic Application of Bidirectional Forwarding Detection (BFD), June 2010.Internet Control Message Protocol (ICMP):        RFC 792—Internet Control Message Protocol, September 1981.        RFC 4443—Internet Control Message Protocol (ICMPv6) for the Internet Protocol Version 6 (IPv6) Specification, March 2006.Network Unreachability Detection (NUD):        RFC 4861—Neighbor Discovery for IP version 6 (IPv6), September 2007.Note: The contents of these eight documents are hereby incorporated by reference herein.Time To Live (TTL) and Hop Limit        \nThere is a field in the IPv4 header, Time To Live (TTL), and in the IPv6 header, Hop Limit, which ensures that a packet caught in a routing loop will eventually be discarded. Plus, when a router 5021, 5022 . . . 502n (for example) receives an IP packet, the router 502k, 5022 . . . 502n decrements the TTL/Hop Limit field. If the result is zero (0), then the packet is discarded.\nProblems\nThe host 500 (for example) which has multiple tunnels 6041, 6042 . . . 604n one to each of its adjunct routers 502k, 5022 . . . 502n must know the operational status of each and every tunnel 6041, 6042 . . . 604n. Furthermore, the tunnels 6041, 6042 . . . 604n which are utilized for communications would depend on the Layer 1/2 active path(s) 5061, 5062 . . . 506n to gain the most efficient routing. This implies that the host's Layer 3 function would need to be cognizant of the Layer 1/2 functions which is not the intent of the ISO OSI Reference Model. This situation can be problematic because some tunnels 6041, 6042 . . . 604n may be intended to be active while others are standby, while some links 5061, 5062 . . . 506n may be intended to be active while other links are standby. And, if the active tunnels 6041, 6042 . . . 604n and the active links 5061, 5062 . . . 506n do not align, then less efficient routing will occur. Referring to FIGS. 8A-8D (PRIOR ART), there are shown two cases—FIGS. 8A and 8B—which lead to efficient routing and two cases—FIGS. 8C and 8D—which lead to less efficient routing. In FIG. 8A, the host 500 and routers 5021 and 5022 have active tunnels 6041 and 6042 and active links 5061 and 5062 which are synchronized resulting in efficient routing of packets. In FIG. 8B, the host 500 and router 5021 have an active tunnel 6041 and an active link 5061 which are synchronized resulting in efficient routing of packets. In FIG. 8C, the host 500 has an active tunnel 6041 and an active link 5061 with router 5021 and the host 500 has an active tunnel 802 but no active link with router 5022 which results in less efficient routing. In FIG. 8C, assuming that one tunnel is all that is needed by the host 500, then the active tunnel 6041 associated with the active link 8001 should be used which avoids the extra hop between the pair of routers 5021 and 5022. In FIG. 8D, the host 500 has an active link 5061 but no active tunnel with router 5021 and at the same time the host 500 has an active tunnel 802 but no active link with router 5022 which results in less efficient routing. Accordingly, there is a need to address these shortcomings and other shortcoming associated with the state-of-the-art."}
{"text": "1. Field of the Invention\nThe present invention relates to a multipolar electrical connector which is to be disposed on a memory card such as an IC card, in order to connect various devices mounted on the memory card with electronic circuits of a computer system.\n2. Description of the Prior Art\nSuch a multipolar electrical connector for a memory card (hereinafter, referred to as merely \"connector\") is attached to the front end portion of the body of the memory card so as to be united with the body of the memory card. The connector has an electrically insulating body wherein a number of pin terminal insertion spaces into which pin terminals of a counter multipolar electrical connector (hereinafter, referred to as merely \"counter connector\") are to be respectively inserted are formed. In each of the pin terminal insertion spaces, a contact piece which functions as a contact and which corresponds to a pin terminal of the counter connector is housed.\nConventionally, connectors having such a configuration, are known, such as those shown in FIGS. 7 to 10. These have a number of contacts (number of contacting points) for pin terminals of a counter connector, and a positional relationship in the front-to-back direction between a contact with which a pin terminal of a counter connector contacts and the rocking fulcrum of a contact piece.\nThe connector shown in FIG. 7 comprises an insulating body A in which pin terminal insertion spaces B are formed. In each of the pin terminal insertion spaces B, a pair of plate-like elongated contact pieces C, C which extend in the pin terminal insertion and extraction direction X are symmetrically arranged. The contact pieces C, C are electrically connected with each other by a short-circuit portion D made of a thin metal plate. The short-circuit portion D and the paired contact pieces C, C are formed into one body by punching or bending a thin metal plate which is common to them.\nThe connector of FIG. 8 comprises a pair of contact pieces C, C and a short-circuit portion D which also are formed by punching a thin metal plate. The edge faces (which are formed as a result of a punching process of the thin metal plate) of the contact pieces C, C function as contacts F, F.\nIn the connector of FIG. 9, contacts F, F of paired contact pieces C, C are positionally shifted in the front-to-back direction, and the base portions of the contact pieces C, C are electrically connected with each other by a short-circuit portion D.\nIn the pair of contact pieces C, C used in the connectors of FIGS. 7, 8 and 9, the contacts F, F corresponding to a pin terminal G1 or G2 are at positions (given positions in the pin terminal extraction direction X1) which are located more forward than rocking fulcrums E, E respectively located in the base portions.\nThe connector shown in FIG. 10 is similar to that of FIG. 7 in that a pair of elongated plate-like contact pieces C, C extend in the pin terminal insertion and extraction direction X and are symmetrically arranged in each of the pin terminal insertion spaces B formed in the insulating body A, in that the contact pieces C, C are electrically connected with each other by a short-circuit portion D made of a thin metal plate, and in that the short-circuit portion D and the contact pieces C, C are formed into one body by punching or bending a thin metal plate which is common to them. In the connector shown in FIG. 10, however, contacts F, F corresponding to a pin terminal G1 or G2 are at positions (given positions in the pin terminal insertion direction X2) which are located more rearward than rocking fulcrums E, E respectively located in the base portions.\nIn FIGS. 7 and 10, B1 designates the body of the counter multipolar electrical connector.\nA number of pin terminals of a counter connector disposed on a computer apparatus in which a memory card is to be used project from the connector body by lengths varying depending on the kind of signals to be processed. This is because, when the pin terminals of the counter connector are inserted into the pin terminal insertion spaces of a connector to connect the counter connector with the connector, the contact timings between the pin terminals through which signals of different kinds pass and the corresponding contact pieces of the connector are temporarily shifted from each other in accordance with the kinds of the signals.\nThe connector described in conjunction with FIG. 7 has higher reliability with respect to that, when the two kinds of pin terminals, the short pin terminal G1 and the long pin terminal G2 indicated by phantom lines in the figure are to be inserted, the timings of contacting the pin terminals G1 and G2 with the contacts F . . . of the corresponding contact pieces C . . . (the contact pieces corresponding to the long pin terminal G2 are not shown) are surely temporarily shifted from each other. More specifically, the pin terminal G1 collides with the contacts F, F immediately before the pin terminal G1 is inserted into the lap position of the paired contact pieces C, C (the position where the contact pieces C, C overlap with each other in the thickness direction T of the body A). It is impossible for the pin terminal G1 to contact any portion of the contact pieces C, C before the pin terminal G1 collides with the contacts F, F. This is applicable also to the long pin terminal G2. As a result, it is possible to assure the above-mentioned temporal shift between the contact timings.\nHowever, the contacts F, F are located more forward than the rocking fulcrums E, E which are respectively located in the base portions of the contact pieces C, C. Accordingly, when, in the insertion process of the pin terminals, the short pin terminal G1 collides with the contacts F, F, for example, a part of the force generated by the collision is received by the rocking fulcrums E, E, whereby a large force is generated so as to impede the insertion of the pin terminal G1. This causes the force (insertion force) required for inserting the pin terminal G1 to be increased. Also for the long pin terminal G2, there arises the same situation as that of the short pin terminal G1. Since the pin terminals G1 and G2 respectively correspond to pairs of the contact pieces C, C and a number of the pin terminals G1 . . . and G2 . . . strike the corresponding contacts F . . . in the insertion process of a multipolar electrical connector, it is necessary to exert a large insertion force, whereby the operability in the insertion process is impaired. To comply with this, a configuration may be adopted in which a single contact piece is allocated to one pin terminal so that the operability in the insertion process is improved. In this configuration, however, each pin terminal contacts with only one contacting point, and therefore contact stability is lowered.\nThe connector described in conjunction with FIG. 8 has a configuration in which the contact pieces C, C are difficult to bend in the directions separating from each other. In this connector, therefore, the force of inserting the pin terminals may be larger than that in the connector of FIG. 7, so that operability in the insertion process is further impaired as compared with that of FIG. 7. Since the edge faces which are formed as a result of a punching process of a thin metal plate function as the contacts F, F, moreover, the maximum width of the contacts F, F is restricted to the thickness of the thin metal plate. When a pin terminal inserted between the contacts F, F is deviated with a somewhat large degree in the width direction of the contacts F, F (the thickness direction, or the direction perpendicular to the sheet in FIG. 8), therefore, it is not possible to absorb the deviation.\nIn the configuration of the connector described in conjunction with FIG. 9, when a pin terminal is inserted, the two contacts F, F sequentially collide with the pin terminal. Therefore, the operability in the insertion process is improved as compared with that of FIG. 7. In order to equalize the effective lengths L1 and L2 of the contact pieces C, C to each other while keeping the contact pressures against the pin terminal equal to each other, however, the contact piece C having the contact F which is located rearward must be shifted backward from the position of the contact piece C having the contact F which is located forward, by the distance corresponding to the shifting distance between the contacts F, F. This increases the whole length of the connector, thereby impeding the miniaturization.\nIn the connector of FIG. 10, the contacts F, F are located more rearward than the rocking fulcrums E, E which are located in the base portions of the contact pieces C, C. According to this configuration, when the short pin terminal G1 collides with the contacts F, F in the insertion process, for example, it occurs only that the paired contact pieces C, C swing about the respective rocking fulcrums E, E, and the situation that the force generated by the collision of the pin terminal G1 against the contacts F, F is received by the rocking fulcrums E, E does not occur. Accordingly, a large force of hindering the pin terminal from being inserted is not generated, and therefore the insertion force for the pin terminal G1 can be smaller than that required in the connectors of FIGS. 7 and 8. Also for the long pin terminal G2, there arises the same situation as that of the short pin terminal G1. Although the pin terminals G1 and G2 respectively correspond to pairs of the contact pieces C, C and a number of the pin terminals G1 . . . and G2 . . . strike the corresponding contacts F . . . in the insertion process of a multipolar electrical connector, therefore, the required insertion force is not so large and operability is hardly impaired.\nHowever, the connector has the following drawback. When the two kinds of pin terminals, the short pin terminal G1 and the long pin terminal G2 indicated by phantom lines in FIG. 10 are to be inserted, it is difficult to assure the temporary shift between the timings of contacting the pin terminals G1 and G2 with the contacts F of the corresponding contact pieces C . . . (the contact pieces corresponding to the long pin terminal G2 are not shown). More specifically, the pin terminal G1 collides with the contacts F, F after the pin terminal G1 is inserted into the lap position of the paired contact pieces C, C. If the pin terminal G1 is bent, therefore, there may arise a situation whereby the short pin terminal G1 contacts the corresponding contact piece C before the long pin terminal G2 collides with the corresponding contacts F, F. When such a situation occurs, it is impossible to assure the temporary shift between the contact timings of contacting the pin terminals G1 and G2 of the two kinds with the contacts.\nAs described above, in a prior art multipolar electrical connector for a memory card, the temporary shift between the timings of contacting contacts with pin terminals having different lengths can be assured at the sacrifice of operability in the insertion process. Further, a prior art multipolar electrical connector has a problem in that the prevention of an impaired operability in the insertion process is liable to be conducted at the sacrifice of the reliability with respect to the assurance of the temporary shift between the contact timings. Moreover, a prior art multipolar electrical connector has another problem in that, when the assurance of the temporary shift between the contact timings is attempted without sacrificing operability in the insertion process, miniaturization is impeded or the reliability of the contacting states between contacts and pin terminals is lowered."}
{"text": "A need exists for a high speed screening device that can be installed in a new and/or existing facility and efficiently screen materials to a tight specification at an industrial minerals processing facilities.\nA further need exists for a screening device that can make multiple gradation cuts, of particulate simultaneously while creating very little dust in the facility.\nA need exists for a device that will prevent human harm during the screening of particulate.\nThe present embodiments meet these needs.\nThe present embodiments are detailed below with reference to the listed Figures."}
{"text": "A current trend in the communications industry is to equip portable electronic devices with a wireless communication interface to communicate in an automated fashion with other devices within its communication range. Typically, technologies such as Bluetooth IEEE 802.15, IEEE 802.11a, IEEE 802.11b, or the like are used and these technologies are fast becoming industry standards. To facilitate universal communication between these devices, universal communication network connectivity arrangements have been developed such as Universal Plug and Play (UPnP). UPnP provides an industry standard architecture for pervasive peer-to-peer network connectivity of electronic devices to allow seamless proximity networking. Data may thus be communicated or exchanged between different devices in an uncontrolled and automated fashion."}
{"text": "In recent years personal computers have become very popular. Now, digital cameras that enable photographed picture information to be easily processed are rapidly becoming popular.\nMore compact and inexpensive digital cameras are being demanded; therefore, it is especially urgent to provide a more compact and inexpensive photographic lens for such cameras. To respond to these requirements, the inventor of the present application has developed a single focus lens that is formed of only three lens elements, as disclosed in Japanese Patent Application 2001-156653. This single focus lens is formed of, in order from the object side, a first lens element having at least one surface aspheric and which has a weak refractive power, a second lens element of positive refractive power and having its image side surface convex, and a third lens element having at least one surface aspheric and which has a weak refractive power. Although the above-mentioned prior art can meet the requirement of being compact and inexpensive, additional improvements are required from the standpoint of aberration correction in order to provide a higher-quality model of single focus lens. Furthermore, as compared with a conventional compact camera that uses film, the design of a lens for a digital camera is much different not only because of the electronic image sensor that receives the image in a digital camera being much smaller, but also because it is necessary in digital cameras to have the light rays be incident onto the image sensor roughly perpendicular to the detecting surface of the image sensor. Further, a high aperture efficacy and a long back focus are required in digital cameras."}
{"text": "Certain printers are capable of receiving communications and printable content via the internet without being connected to a desktop computer, notebook computer, or other host computing device. An advantage of such a printer print jobs can be received for printing from other computing devices located anywhere around the globe.\nThe same part numbers designate the same or similar parts throughout the figures."}
{"text": "The invention concerns a method and apparatus for rapidly solidifying and cooling molten metal oxides and metal oxide-base compositions by continuous casting.\nIn many industrial processes an attempt is made to obtain rapid solidification and cooling of melted products based on metal oxides, for reasons either of safety or of convenience: for example, when exhausted drosses or slags from electro-metallurgical operations have to be taken rapidly to the tip without the risk of transporting them in the liquid or partially solidified state, or for structural reasons, when rapid cooling is necessary to give the product certain physical properties. For example, abrasives based on corundum and on corundum and zirconia are known to have mechanical properties which improve commensurately with the rapidity with which the molten product, manufactured by fusion in an electric furnace, is cooled. This fact has been pointed out particularly in French Pat. No. 1 332 975 in the name of CARBORUNDUM Company, which recommends casting in a graphite mold, in French Pat. No. 2 127 231 in the name of PECHINEY directed to casting the melted product onto steel balls, in French Pat. No. 2 133 595 in the name of TREIBACHER CHEMISCHE WERKE AG directed to casting onto iron spheres, and in French Pat. No. 2 166 082 in the name of NORTON directed to casting between solid metal plates 1.6 to 12.7 mm apart.\nFor obvious reasons it would be desirable for the rapid cooling to take place continuously. In the case of drosses and slags, as a matter of fact, cooling in ingot molds, which is the usual practice, is very slow and necessitates a large number of molds, which take up space in the workshop. In the case of abrasives with fine crystallization, casting between closely spaced metal plates or onto steel balls leads to very low productivity.\nA certain number of processes for cooling rapidly by continuous casting have previously been described.\nFrench Pat. No. 1 319 102 to NORTON describes the casting of electro-melted alumina, in a thin layer, onto a hollow cylinder which is cooled by sprinkling water onto its internal surface.\nU.S. Pat. No. 3,993,119, and French Pat. No. 2 290 266 to NORTON, describe an arrangement for casting between stationary, thick metal plates, through their base onto an endless conveyor. The conveyor passes over rollers which causes the plates to move apart and the blocks of cooled abrasive to be separated and drop into a receiving means.\nFrench Pat. No. 2 422 462 to NORTON describes another arrangement for continuous casting between cooled metal belts which are kept in contact with the solidifying abrasive by a series of guiding rollers.\nThese various pieces of apparatus are mechanically complex and ill adapted to cheap, continuous mass production of electro-melted abrasives with fine crystallization."}
{"text": "Wireless communication systems are widely deployed to provide various types of communication content such as voice, data, and so on. These systems may be multiple-access systems capable of supporting communication with multiple users by sharing the available system resources (e.g., bandwidth and transmit power). Examples of such multiple-access systems include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, and orthogonal frequency division multiple access (OFDMA) systems.\nGenerally, a wireless multiple-access communication system can simultaneously support communication for multiple wireless terminals. Each terminal communicates with one or more base stations via transmissions on the forward and reverse links. The forward link (or downlink) refers to the communication link from the base stations to the terminals, and the reverse link (or uplink) refers to the communication link from the terminals to the base stations. This communication link may be established via a single-in-single-out, multiple-in-signal-out or a multiple-in-multiple-out (MIMO) system.\nUniversal Mobile Telecommunications System (UMTS) is one of the third-generation (3G) cell phone technologies. UTRAN, short for UMTS Terrestrial Radio Access Network, is a collective term for the Node-B's and Radio Network Controllers which make up the UMTS radio access network. This communications network can carry many traffic types from real-time Circuit Switched to IP based Packet Switched. The UTRAN allows connectivity between the UE (user equipment) and the core network. The UTRAN contains the base stations, which are called Node Bs, and Radio Network Controllers (RNC). The RNC provides control functionalities for one or more Node Bs. A Node B and an RNC can be the same device, although typical implementations have a separate RNC located in a central office serving multiple Node B's. Despite the fact that they do not have to be physically separated, there is a logical interface between them known as the Iub. The RNC and its corresponding Node Bs are called the Radio Network Subsystem (RNS). There can be more than one RNS present in an UTRAN.\n3GPP LTE (Long Term Evolution) is the name given to a project within the Third Generation Partnership Project (3GPP) to improve the UMTS mobile phone standard to cope with future requirements. Goals include improving efficiency, lowering costs, improving services, making use of new spectrum opportunities, and better integration with other open standards. The LTE system is described in the Evolved UTRA (EUTRA) and Evolved UTRAN (EUTRAN) series of specifications.\nSemi-persistent scheduling (SPS) is a set of techniques for efficiently assigning resources for periodic traffic in a wireless communication system to support resource assignment with as little overhead as possible in order to improve system capacity.\nMobile Internet Protocol (Mobile IP or MIP) is a communication protocol that enables transparent routing of data packets to mobile devices in a wireless communication system. Under the Mobile IP protocol, a device can register with a home agent (HA), through which the device obtains a “home” IP address. The home address of the device can then be utilized to route data packets to and/or from the device regardless of the location of the device within a wireless communication network. Conventionally, a mobile device can register with a HA by first discovering a global IP address of the HA and subsequently setting up a security association with the HA based on its discovered IP address. Upon association with the HA, the device can signal updates to the HA relating to the location and/or status of the device. These updates can be utilized by the HA to provide data packets to the device, either directly or indirectly via an access point of a disparate network to which the device has moved.\nHowever, complications arise regarding Mobile IP and policy charging control (PCC) interactions. In particular, the entities involved in the mobility management (i.e., mobile node and home agent) are different from the entities involved in quality of service (QoS) and policy (i.e., PCC) management, specifically named PCEF (Policy and Charging Enforcement Function) and PCRF (Policy and Charging Rules Function) respectively. When Mobile IP is not used, there is only one IP address used by the mobile node and this is communicated to the QoS management entities (i.e. PCRF). Based on this address the PCRF establishes the correct QoS in the system. However, when a mobile node performs a MIP registration with the HA, one of the IP address of the mobile node can change and the QoS management entities (i.e., PCRF) are not aware of that as they are not involved. Based on this situation, the PCRF does not know which addresses have been registered by the mobile node. Therefore, the PCRF cannot determine which PCC rules (e.g., IP flows and related policies) need to be sent to the PCEF located at the access gateway for a respective MIP registration."}
{"text": "As computers and data processing equipment have grown in capability, users have developed applications that place increasing demands on the equipment. Thus, there is a continually increasing need to process more information in a given amount of time. One way to process more information in a given amount of time is to process each element of information in a shorter amount of time. As that amount of time is shortened, it approaches the physical speed limits that govern the communication of electronic signals. While it would be ideal to be able to move electronic representations of information with no delay, such delay is unavoidable. In fact, not only is the delay unavoidable, but, since the amount of delay is a function of distance, the delay varies according to the relative locations of the devices in communication.\nSince there are limits to the capabilities of a single electronic device, it is often desirable to combine many devices, such as memory components, to function together to increase the overall capacity of a system. However, since the devices cannot all exist at the same point in space simultaneously, consideration must be given to operation of the system with the devices located diversely over some area.\nTraditionally, the timing of the devices' operation was not accelerated to the point where the variation of the location of the devices was problematic to their operation. However, as performance demands have increased, traditional timing paradigms have imposed barriers to progress.\nOne example of an existing memory system uses DDR (double data rate) memory components. The memory system includes a memory controller and a memory module. A propagation delay occurs along an address bus between the memory controller and the memory module. Another propagation delay occurs along the data bus between the memory controller and the memory module.\nThe distribution of the control signals and a control clock signal in the memory module is subject to strict constraints. Typically, the control wires are routed so there is an equal length to each memory component. A “star” or “binary tree” topology is typically used, where each spoke of the star or each branch of the binary tree is of equal length. The intent is to eliminate any variation of the timing of the control signals and the control clock signal between different memory components of a memory module, but the balancing of the length of the wires to each memory component compromises system performance (some paths are longer than they need to be). Moreover, the need to route wires to provide equal lengths limits the number of memory components and complicates their connections.\nIn such DDR systems, a data strobe signal is used to control timing of both data read and data write operations. The data strobe signal is not a periodic timing signal, but is instead only asserted when data is being transferred. The timing signal for the control signals is a periodic clock. The data strobe signal for the write data is aligned to the clock for the control signals. The strobe for the read data is delayed by delay relative to the control clock equal to the propagation delay along the address bus plus the propagation delay along the data bus. A pause in signaling must be provided when a read transfer is followed by a write transfer to prevent interference along various signal lines used. Such a pause reduces system performance.\nSuch a system is constrained in several ways. First, because the control wires have a star topology or a binary tree routing, reflections occur at the stubs (at the ends of the spokes or branches). The reflections increase the settling time of the signals and limit the transfer bandwidth of the control wires. Consequently, the time interval during which a piece of information is driven on a control wire will be longer than the time it takes a signal wavefront to propagate from one end of the control wire to the other. Additionally, as more modules are added to the system, more wire stubs are added to each conductor of the data bus, thereby adding reflections from the stubs. This increases the settling time of the signals and further limits the transfer bandwidth of the data bus.\nAlso, because there is a constraint on the relationship between the propagation delays along the address bus and the data bus in this system, it is hard to increase the operating frequency without violating a timing parameter of the memory component. If a clock signal is independent of another clock signal, those clock signals and components to which they relate are considered to be in different clock domains. Within a memory component, the write data receiver is operating in a different clock domain from the rest of the logic of the memory component, and the domain crossing circuitry will only accommodate a limited amount of skew between these two domains. Increasing the signaling rate of data will reduce this skew parameter (when measured in time units) and increase the chance that a routing mismatch between data and control wires on the board will create a timing violation.\nAlso, most DDR systems have strict limits on how large the address bus and data bus propagation delays may be (in time units). These are limits imposed by the memory controller and the logic that is typically included for crossing from the controller's read data receiver clock domain into the clock domain used by the rest of the controller. There is also usually a limit (expressed in clock cycles) on how large the sum of these propagation delays can be. If the motherboard layout makes this sum too large (when measured in time units), the signal rate of the system may have to be lowered, thereby decreasing performance.\nIn another example of an existing memory system, the control wires and data bus are connected to a memory controller and are routed together past memory components on each memory module. One clock is used to control the timing of write data and control signals, while another clock is used to control the timing of read data. The two clocks are aligned at the memory controller. Unlike the previous prior art example, these two timing signals are carried on separate wires.\nIn such an alternate system, several sets of control wires and a data bus may be used to intercouple the memory controller to one or more of the memory components. The need for separate sets of control wires introduces additional cost and complexity, which is undesireable. Also, if a large capacity memory system is needed, the number of memory components on each data bus will be relatively large. This will tend to limit the maximum signal rate on the data bus, thereby limiting performance.\nThus, a technique is needed to coordinate memory operations among diversely-located memory components."}
{"text": "Total or near-total occlusions in arteries can prevent all or nearly all of the blood flow through the affected arteries. It has been estimated that 5% to 15% of patients on whom percutaneous coronary angioplasty (PTCA) is attempted are found to have chronic total occlusions (CTO's) of at least one coronary artery. Chevalier, B. et al., Chronic Total Occlusion, The Paris Course on Revascularization, Pages 131–148 (May 2000) In patients who suffer from coronary CTO's, the successful performance of a PTCA is a technical challenge. The factor that is most determinative of whether the interventionalist can successfully perform PTCA on patient who presents with a coronary CTO is the interventionalist's ability (or inability) to advance a suitable guidewire from a position within the lumen of the artery upstream of the lesion, across the lesion (i.e., either through the lesion or around it), and then back into the artery lumen at a location downstream of the lesion.\nIn some instances, such as where the occlusive matter is soft or where the occlusion is less than total, the guidewire can simply be pushed through the occlusive matter itself, thereby allowing the guidewire to remain within the artery lumen. However, in other cases, such as where the artery is totally occluded by hard, calcified atherosclerotic plaque, the guidewire may tend to deviate to one side and penetrate through the intima of the artery, thereby creating a neo-lumen through the sub-intimal space (i.e., within the wall of the artery between the intima and adventitia). In these cases, after the distal end of the guidewire has been advanced to a position distal to the lesion, it is then necessary to divert or steer the guidewire from the sub-intimal space back into the lumen of the artery at a location downstream of the lesion. This process of causing the guidewire to reenter the artery lumen is often difficult and others have proposed various means for dealing with such problem. For example, PCT International Publication No. WO 00/18323 describe techniques where a catheter that has a through lumen terminating distally in a laterally directed outlet port is advanced into the sub-intimal space past the lesion and a penetrator or guidewire is then advanced through the catheter, out of the laterally directed outlet port, and back into the lumen of the artery. WO 00/18323 further states that optionally a wire may be passed through a lumen of the catheter and that such wire may further comprise an imaging apparatus such as an ultrasonic imaging means. Also, WO 00/18323 states that the catheter may include a marker near its distal end that is visible on fluoroscopy and that such marker may be configured to permit visual determination of the rotational orientation of the distal end of the catheter when viewed as a two-dimensional fluoroscopic image. However WO 00/18323 does not describe any means for correlating the rotational orientation of the catheter to the trajectory on which the penetrator or guidewire will subsequently advance. A number of variable can affect the trajectory on which the penetrator or guidewire will advance, including the location of the outlet aperture from which the penetrator or guidewire exits the catheter body, the density and resistance of the tissue through which the penetrator or guidewire advances and any bias, shape or curvature that the penetrator or guidewire will assume after it is no longer constrained by the catheter.\nAnother approach to bypassing a CTO, as described in U.S. Pat. No. 6,068,638 (Makower), U.S. Pat. No. 5,830,222 (Makower) and U.S. Pat. No. 6,190,353B1 (Makower, et al.) utilizes a tissue penetrating catheter device that has an on board guidance element that allows a tubular penetrator cannula to penetrate outwardly through the intima of the artery wall upstream of the lesion after which a member is advanced trough the cannula and bores through tissue located outside of the vessel's intima, to a position distal to the obstructive lesion, and then reenters the lumen of the artery at a location downstream of the lesion. Tissue penetrating catheters of this type, having on-board guidance and/or orientation capabilities and penetrator cannulas that are advanceable laterally from the catheter body, have previously been devised. (Transaccess® catheters by Transvascular,® Inc., Menlo Park, Calif.)\nThere exists a need in the art for the development of new and better techniques for catheter-based treatment of CTO and other total or near-total obstruction of arteries or other anatomical conduits, especially with respect to improving the available techniques for causing a guidewire to reenter the lumen of an artery from a location within the sub-intimal space or even outside of the artery wall and for facilitating verification of the intended re-entry into the artery lumen and prompt delivery of arterial blood or some other oxygenated perfusate to tissues that would normally receive arterial blood flow from the obstructed artery."}
{"text": "A system (hereinafter, referred to as a “wireless power feeding system”) using wireless power transmission in which power is supplied to an electric apparatus in a wireless manner without using a power cord or the like is being put into practical use. For example, an electromagnetic induction type wireless power feeding system using electromagnetic induction between antennae (coils) which are disposed to be separated from each other, or a magnetic resonance type wireless power feeding system using electromagnetic field resonance coupling is known. In addition, near field communication (NFC) is known as a standard regarding a wireless communication technique for transmitting information in a wireless manner, and IC cards or small portable terminal apparatuses conforming to the NFC standard also start to being widespread.\nThe magnetic resonance type wireless power feeding system is implemented by using a resonance circuit including coils and capacitors. The magnetic resonance type wireless power feeding system can have a longer transmission distance between a power transmission coil and a power reception coil than in an electromagnetic induction type of the related art by increasing a Q value of the coil, and has a feature of being strong to positional deviation of the power transmission coil and the power reception coil.\nIn addition, the magnetic resonance type wireless power feeding system has a feature of being hardly influenced by a foreign substance present between a power transmission side and a power reception side. However, since an amount of transmission power absorbed by a foreign substance is increased depending on the kind of foreign substance, power transmission efficiency is reduced, and there is also a concern that the foreign substance may be damaged due to heat generation caused by the power absorbed by the foreign substance.\nThe following PTLs 1 and 2 disclose techniques for detecting a foreign substance or an abnormality in a power feeding target in the magnetic resonance type wireless power feeding system.\nFor example, PTL 1 discloses a technique for detecting a foreign substance present between a power transmission side and a power reception side in a magnetic resonance type wireless power feeding system in which a power transmission coil is provided on a road, and a power reception coil is mounted on a vehicle, so that the vehicle is charged in a wireless manner. Specifically, according to the configuration disclosed in PTL 1, an impedance value is estimated by detecting a terminal voltage of a battery which is mounted in the vehicle and is charged on the basis of received power, and a notification of the presence of a foreign substance is sent and power feeding is also stopped in a case where a difference between the estimated value and an actually measured impedance value exceeds a threshold value.\nIn addition, PTL 2 discloses a technique for detecting an abnormality in a vehicle by detecting an increase in reflected power in a magnetic resonance type wireless power feeding system in which a power transmission coil is provided on a road, and a power reception coil is mounted on the vehicle, so that the vehicle is charged in a wireless manner. Specifically, according to the configuration disclosed in PTL 2, if a suspicious person gets in the vehicle which is being charged, impedance mismatching occurs due to a change in a vehicle height, thus reflected power increases, and a notification of an abnormality in the vehicle is sent when the increase in the reflected power is detected."}
{"text": "This invention relates to wet-spun plural-end elastane yarn having a yarn linear density of up to 2500 dtex, high surface lustre and a ribbony cross section and formed from bonded-together individual filaments having an oval to circularly round cross section, the width of the elastane yarn in cross section being at least four times the thickness of the elastane yarn.\nElastane yams are synthetic filament yarns which are produced as mono-or multifilamentsxe2x80x94dependent on the intended usexe2x80x94within the linear density range from. 11 to about 2500 dtex. Most elastane yams are coalesced multifilament yarns and possess in their grooved or fluted surface a very specific textile character (compare P. A. Koch xe2x80x9cElastanfasernxe2x80x9d in Chemiefasern/Textilindustrie February 1979, page 97-98).\nThe cross-sectional shape of elastane yarn is substantially dependent on the production process. In the case of dry-spun elastane yarn, the cross-sectional shape will alter in some instances appreciably with the yarn linear density, relatively fine linear densities usually having roundish cross sections and coarser linear densities from about 500 dtex usually having oval to dumbbell-shaped cross sections. This is believed to be due to differences in the degree of evaporation of the spinning solution solvent under the particular spinning conditions. The particular spinning conditions thus influence the onset of coagulation and hence also the degree of shrinkage in the yam\"\"s cross section.\nThese correlations between the cross-sectional shape and the solvent evaporation rate for dry spinning from polymer solutions, for example of acrylic fibres, are extensively described in the OPI document EP-A-31 078.\nUsing dry spinning for elastane yarn for medium and coarse titers of about 160 dtex, utilizing jets with several jet openings for glued multifilaments, will usually provide roundish to oval overall cross-sectional shapes, averaged over the entire cross section, medium and higher linear densities produced by wet spinning will usually prosses ribbony cross sections in multiple rows, strung up like pearls on a string.\nThe ribbony cross sections can be advantageous and disadvantageous in the desired final article. The ribbony shape is disadvantageous especially in wide fabric, for example articles having a large area, when its sparkle is unwelcome. On the other hand, ribbony cross sections are also of advantage, for example for the production of diapers. Owing to the wide cross-sectional shape, it is particularly readily possible to apply adhesives to the yarn which lead to firm incorporation and adhesion of the elastic yarn in the diaper article.\nWet-spun elastanes generally have an as-spun filament linear density which is between 15 and 25 dtex.\nKnown elastane yarn, in addition, has a comparatively dull surface due to the superficial and cross-sectional structure of its individual filaments. In certain applications for the production of textile sheet material, however, particular importance attaches to the visual quality of the fibres, especially their lustre. An objective measure of the lustre is the so-called lustre number which is defined hereinbelow. The lustre number of existing wet-spinnable elastane yarn is at best on the order of 10 to 20.\nThe invention has for its object to provide elastane yarns having a yarn linear density of up to 2500 dtex and a process for their production which have similar mechanical properties to known elastane yarns, a ribbony cross section and a very high surface lustre.\nIt has now been surprisingly found that raising the filament count whilst keeping the final yarn linear density constant leads to a distinct broadening of the ribbon shape of wet-spun plural-end elastane yarn by up to about 40% compared with the width of known elastane yarn when the as-spun filament linear density is less than 15 dtex, especially not more than 10 dtex.\nThis object is achieved according to the invention by wet-spun plural-end elastane yarn having a yarn linear density of up to 2500 dtex, high surface lustre and a ribbony cross section and formed from bonded-together individual filaments having an oval to circularly round cross section, the width of the elastane yarn in cross section being at least four times the thickness of the elastane yarn and the as-spun filament linear density being less than 15 dtex.\nIn particularly suitable elastane yarn, the width of the elastane yarn in cross section is at least five times, preferably at least eight times, particularly preferably at least ten times, the thickness of the elastane yarn.\nIn elastane yarn of particularly uniform surface quality, the as-spun filament linear density is not more than 10 dtex, preferably not more than 5 dtex.\nPreference is given to elastane yarn which has a lustre number of at least 20, preferably of least 40, particularly preferably of at least 100.\nPreferably, the elastane yarn comprises polyurethane comprising not less than 85% by weight of segmented polyurethane.\nSegmented polyurethane is in particular a segmented polyurethane based on polyethers, polyesters, polyetheresters, polycarbonates or mixtures thereof.\nThe invention further provides a wet spinning process for producing plural-end elastane yarn having a yarn linear density of up to 2500 dtex, high surface lustre and a ribbony cross section by spinning and especially at least 25% strength by weight elastane solution into a coagulation bath, washing and optionally drawing the yarn formed, drying, setting, spin finishing and winding the yarn, characterized in that the as-spun filament linear density is less than 15 dtex and the individual filaments are converged on a diverting roller in the coagulation bath by means of a comb-type thread guide.\nIn the preferred process, the as-spun filament linear density is not more than 10 dtex, preferably not more than 5 dtex, particularly preferably not more than 3 dtex so to obtain an almost circularly round cross-sectional shape for the individual filaments.\nIt is particularly advantageous when the elastane yarn is additionally pressed by a press roller onto the diverting roller after the converging on the diverting roller in the coagulation bath.\nThe coagulation bath press roller contact pressure is preferably within the range from 1 to 5 bar (0.1 to 0.5 MPa), in particular for elastane yam having a linear density of up to 1000 dtex 2 to 3 bar (0.2 to 0.3 MPa).\nThe choice of the as-spun filament linear density, as well as broadening the ribbon shape of the elastane yarn, surprisingly also modifies the cross section of the individual filaments. Whereas an as-spun filament linear density of more than 15 dtex is observed to produce, for example, kidney-shaped filament cross sections, an as-spun filament linear density below 15 dtex, especially of about 10 dtex or less, is found to be accompanied by a change in the cross-sectional shape of the individual filaments into an oval to circularly round shape. At an as-spun filament linear density of less than about 5 dtex, virtually all the filaments present have a round cross-sectional shape. FIGS. 1a to 1c show the increase in the ribbon width with an increasing filament count for the 800 dtex yarn under otherwise identical spinning conditions. Magnification: 1 cm=about 70 micron. FIGS. 2a to 2c show individual filament cross sections for various as-spun filament linear densities for the 800 dtex yarn. Magnification: 1 cm=30 micron. As is discernible from FIGS. 2a to 2c, wet spinning an 800 dtex elastane yarn from a 385 hole jet using an as-spun filament linear density of 4.2 dtex gives round filament cross sections. The surface of such elastane yarn is smoother, has a lower groove depth and hence possesses a particularly high lustre. While an 800 dtex elastane yarn spun from a 172 hole jet using an as-spun filament linear density of 9.4 dtex possesses a lustre number of 60, the lustre number increases to 100 for an as-spun filament linear density of 4.1 dtex (cf. Tab. I1, Examples 6 and 8).\nThe distinctly higher lustre of the elastane yarn of the invention, which is particularly noticeable in sheetlike articles, is attributable to the smooth and unfibrillated surface of the filaments. The elastane yarn, moreover, has only weakly developed flutes parallel to the fibre axis, which are not interrupted, so that light incident on the elastane yarn is reflected by the individual filaments in a directional manner.\nReducing the as-spun filament linear density whilst increasing the filament count in spinning causes as stated a broadening of the elastane yarn. In the same way, applying additional pressure on the yarn, after coagulation in the coagulation bath, by means of a special squeeze roller which rests on a diverting roller for the yams which is disposed in the coagulation bath leads to a further increase in the ribbon width of about 30-40% depending on the as-spun filament linear density (cf. Table 2, Examples 6-9).\nIn general, a yarn linear density of up to 2500 dtex requires a contact pressure of 1 to 5 bar (0.1 to 0.5 MPa) so that this effect may be obtained. In the case of a yam linear density of up to 1000 dtex, a contact pressure of 2 to 3 bar (0.2 to 0.3 MPa) has proved particularly suitable. A particularly suitable means has proven to be an elastomeric press roller of about 50-100 Shore hardness, whose contact pressure is pneumatically controllable via compressed air.\nThe invention further provides for the use of the wet-spun elastane yarn according to the invention for producing textile goods, preferably sheetlike goods, especially for producing consecutive course formation knits, wovens or synchronous course formation knits.\nThe invention also provides for the use of the wet-spun elastane yarn according to the invention for producing disposable hygiene articles, especially for producing diapers.\nThe ribbon width is reported in the Examples in xe2x80x9cmicronxe2x80x9d (xcexcm). 1 cm is in each case equal to 10,000 micron.\nThe as-spun filament linear density (ASFLD) is calculated as follows:   ASFLD  =                    F        xc3x97        K        xc3x97        0.94        xc3x97        100                    A        xc3x97        Z              ⁢          (      dtex      )      \nwhere\nF=pumpage rate of spinning solution (ccm/min)\nK=concentration of spinning solution (% by weight)\nA=coagulation of speed (n/min)\nz=number of jet holes\nThe as-spun filament linear density is a measure of the filament weight immediately following exit of the spinning solution into the coagulation bath.\nMeasurement of lustre number\nThe elastane yarns are wound flat onto black metal plates so that the surface of the plate is completely covered by elastane yarn. The luminance of the wound yarn is measured as a function of the scattering angle using a goniphotometer. The angle of incidence is 45xc2x0. A plate of barium sulphate serves as the non-lustrous standard.       Lustre    ⁢          xe2x80x83        ⁢    number    ⁢          xe2x80x83        ⁢    z    =            Reflectance      ⁢              xe2x80x83            ⁢      of      ⁢              xe2x80x83            ⁢      sample      ⁢              xe2x80x83            ⁢      at      ⁢              xe2x80x83            ⁢      45      ⁢              xe2x80x83            ⁢      degrees              Reflectance      ⁢              xe2x80x83            ⁢      of      ⁢              xe2x80x83            ⁢      sample      ⁢              xe2x80x83            ⁢      at      ⁢              xe2x80x83            ⁢      0      ⁢              xe2x80x83            ⁢      degrees      \nThe Examples hereinbelow illustrate the invention. Parts and percentages are by weight, unless otherwise stated."}
{"text": "1. Field of the Invention\nThe present invention relates to a licensee notification system employed for the sale of software using a high speed communication network such as B-ISDN and a large-capacity storage medium such as a CD-ROM.\n2. Description of the Related Art\nWith the development of high speed communication technology such as B-ISDN (broad-band integrated services digital network) and high-capacity storage media such as CD-ROMs (compact disk read only memory) such means can now be used to distribute computer programs or video data or audio data. For example, video works which were previously supplied on video tape are now being sold stored on CD-ROM. Also, game programs etc, which contain a large amount of picture data, are being sold stored on CD-ROM. The same applies to high speed communication networks, in which the software supplier can now distribute the software by various methods. One of these methods of software sales is the so-called \"locked software\" sales system. In the locked software sales system, a CD ROM on which are stored a large number of software items whose functions are restricted is sold cheaply. By using the various items of software on the CD-ROM that is purchased, in a condition with the functional restrictions imposed, the end user is able to make a decision as to whether or not he needs each software item. Then, if the end user does require the software, he obtains (purchases) a restriction-removal code corresponding to this software from a management center operated by the software distributor, and is able to use this restriction-removal code to remove the functional restrictions on the software.\nSuch a sales system may be implemented, as a specific example, using the software sales system shown in FIG. 10. As shown in this Figure, this software sales system comprises user terminals 31 and management center 32. The user terminal 31 and the management center 32 are connected by means of a communication circuit.\nWhen actually purchasing the software (i.e. when purchasing a restriction-removal code), the end user, using a user ID etc, sets up a communication path with the management center and executes the prescribed procedure required to request that a restriction-removal code be sent to the user terminal 31. This procedure includes the input of a \"contents ID\", which is information for identifying the software item that is to be purchased actually. In response to the execution of such a procedure, the user terminal 31 sends to the management center 32 the contents ID and for example the characteristic information of the user, consisting of the ID of the CPU provided in user terminal 31.\nWithin the management center 32, there is provided a software database (software DB) in which software decoding keys employed for encoding the various software items are stored in association with the contents ID. When a contents ID is received from user terminal 31, the software decoding key corresponding to the contents ID is read from software database 33. Also, encoding unit 34 in management center 32 generates a user individual key by encoding the user characteristic information from user terminal 31 by the key \"Ks\". Encoding unit 35 sends the results of the encoding of the software decoding key from software database 33 to user terminal 31 as restriction-removal code, using the user individual key from encoding unit 34.\nEncoding unit 36 in user terminal 31 generates a user individual key by encoding the user characteristic information with the key \"Ks\". Decoding unit 37 uses the user individual key generated by encoding unit 36 to decode the restriction removal code from management center 32, thereby generating the software decoding key. Installation unit 38 then uses this software decoding key to decode the software in CD-ROM corresponding to the contents ID sent to center terminal 32: thus the software is put in a condition where it can be used with the functional restrictions removed, and, in this form, is installed on to a storage device such as a hard disk device.\nWith such a software sales system, it is possible to determine the software item to be purchased after actually ascertaining its contents: thus, the possibility that the purchased software might be completely different from that intended, as could happen if the purchase were made solely on the basis of the details contained in a catalogue, can be completely eliminated. Also, since the software on the CD ROM is stored in a form which is not executable without knowing special information, illicit installation can be prevented.\nHowever, once the software has been installed, it is an extremely easy operation to copy this. Thus, the problem has arisen of unscrupulous persons copying the software without the consent of the software supplier. Various methods (so-called protection methods) of preventing such illicit copying are known but there is no way to prevent illicit copying by a person possessing knowledge at the level of the BIOS (basic input/output system). Whichever method is used, it can do no more than make it more difficult to perform illicit copying.\nFor this reason, software is sold in which the name of the authorized user is displayed on start-up, with the object of preventing illicit copying psychologically rather than physically. That is, the aim is to prevent illicit copying of software by displaying the name of the authorized user of the software when the illicitly copied software is executed.\nHowever, even with such software, if the copying is inclusive of the installation software that sets the user name, when the software is run, it can be made to display the name of the person who made the illicit copy: thus, sufficient effectiveness in preventing illicit copying was not obtained."}
{"text": "Omotically active compounds (osmotics) are used widely in pharmacy and medicine. For example, osmotics are used to adjust the tonicity of drugs, in particular parenteral medications, where the osmotic pressure of a drug is adjusted to be hypotonic, hypertonic or isotonic, depending on how it is administered. For example, the osmotic pressure of a parenteral drug solution may be adjusted to match the osmotic pressure of human blood by adding an osmotic agent (isosmotic solutions).\nFurthermore, osmotics are used in dialysis therapy, in particular in peritoneal dialysis, to withdraw excess water from the dialysis patient.\nThe peritoneal dialysis method is based on the fact that a solution containing osmotically active compounds is introduced into the abdominal cavity of a dialysis patient through a catheter. The solution is left in the patient's abdominal cavity, where it manifests its osmotic effect for a certain period of time (usually a few hours). In other words, endogenous water is withdrawn from the patient into the abdominal cavity. After a certain dwell time, the peritoneal dialysis solution, which is then dilute, is drained out through a catheter.\nThis principle is employed in various methods of peritoneal dialysis therapy. For example, the methods of intermittent peritoneal dialysis (IPD), nocturnal intermittent peritoneal dialysis (NIPD), continuous cyclic peritoneal dialysis (CCPD) or continuous ambulant peritoneal dialysis (CAPD) may be employed as needed. In IPD, NIPD and CCPD, instruments which support the patient in performing the peritoneal dialysis method are used. CAPD is a manual method.\nBy adding osmotically active compounds, it should be ensured in particular that the osmotic pressure of the peritoneal dialysis solution is high enough during the entire dwell time in the abdominal cavity to withdraw water from the patient. In other words, water goes from the patient's circulation into the abdominal cavity (ultrafiltration).\nHowever, because of the transfer of water into the abdominal cavity, dilution of the peritoneal dialysis solution introduced there necessarily occurs. As a result of this dilution, there is a decline in the concentration of the osmotically active compound and thus also the osmotic pressure of this solution.\nIf the osmotic pressure of the peritoneal dialysis solution drops because of this dilution, this in turn results in a decrease in or even a complete standstill of the transfer of water into the abdominal cavity per unit of time. In these cases, there is no longer an effective withdrawal of water with advancing dwell time of the peritoneal dialysis solution in the patient's abdominal cavity.\nThrough absorption of osmotically active compounds into the patient's bloodstream, the direction of the transfer of water may even be reversed, i.e., water goes from the abdominal cavity into the patient's bloodstream (negative ultrafiltration). This is the case when the dilute peritoneal dialysis solution in the abdominal cavity has a lower osmotic pressure than the patient's endogenous water (e.g., blood).\nBy adding suitable osmotically active compounds to the peritoneal dialysis solution, the osmotic pressure may be maintained for a treatment time, which is suitable for peritoneal dialysis, so that there is no excessive decline in ultrafiltration within the dwell time of the solution in the abdominal cavity. A negative ultrafiltration is thus also largely prevented.\nThe solutions used in the peritoneal dialysis therapy usually contain sugar monomers or polymers, for example, glucose or polyglucose (e.g., starch derivatives) as osmotically active compounds.\nU.S. Pat. No. 4,889,634 relates to a peritoneal dialysis solution containing hydroxypropyl-β-cyclodextrin.\nJP 8071146 relates to a peritoneal dialysis solution containing α- or γ-cyclodextrin, 2-hydroxyethyl ether, 2-hydroxypropyl ether, 6-O-α-glucosyl or 6-O-α-maltosyl derivatives of α-, β- and γ-cyclodextrin.\nHowever, the state of the art describes cyclodextrin compounds which are not selectively substituted. Instead these cyclodextrin derivatives are mixtures of a wide variety of compounds (different degrees of substitution, different positional isomers) because the substitution of cyclodextrins takes place selectively on certain hydroxyl groups of the cyclodextrins only to a very limited extent."}
{"text": "1. Field of the Invention\nThe present invention concerns a dielectric ceramic composition mainly comprising barium titanate (BaTiO3) and a multi-layer ceramic capacitor using the same and, particularly, a multi-layer ceramic capacitor having internal electrodes formed of Ni or an Ni alloy and an electrostatic capacity with less temperature dependence.\n2. Description of Related Art\nFor multi-layer ceramic capacitors used for electronic equipment such as portable equipment and telecommunication equipment, demand for reduction of size and increase of the capacitance has increased more and more. Such small-sized multi-layer ceramic capacitors having large capacity includes a multi-layer ceramic capacitor in which the internal electrode is constituted with Ni, for example, as shown in JP-A-2001-39765. Such a multi-layer capacitor can provide a multi-layer capacitor that satisfies an X7R property (satisfying the rate of permittivity change within ±15% within a temperature range of from -55° C. to 125° C. with 25° C. as a reference), having a permittivity of 2500 or higher, and having excellent reliability even when it is lamellated.\nHowever, reliability under severe circumstances has been required in recent years for such multi-layer ceramic capacitors depending on the application use. For example, the multi-layer ceramic capacitors have been used for car-mounted electronic equipment such as electronic engine control units and anti-lock brake systems mounted in engine rooms of cars. In such car-mounted electronic equipment, since a stable operation is demanded under conditions of a low temperature lower than −20° C. or higher than +130° C., good temperature stability even under such conditions has been required also for the multi-layer ceramic capacitors used for the car-mounted electronic equipment.\nFor coping with such a demand, a dielectric ceramic composition and a multi-layer ceramic capacitor satisfying the X8R property (rate of permittivity change within ±15% within a temperature range of from -55° C. to 150° C. with 25° C. as a reference) are proposed in, for example, JP-A-2005-272263.\nHowever, the dielectric ceramic composition and the multi-layer ceramic capacitor described above involved a problem that a sufficient specific resistance can not always be obtained in a case of further reducing the thickness of dielectric ceramics between internal electrodes for further reducing the size and increasing the capacitance. Particularly, they involved a problem that the specific resistance is not sufficient under a high temperature circumstance exceeding 125° C."}
{"text": "1. Field of the Invention\nThis invention relates to airbrasive devices, and more particularly to devices for containing abrasive materials expelled by a gas-abrasive apparatus. The invention is particularly useful for dental applications.\n2. Background of the Invention\nThe use of sandblasting devices to contact various surfaces has been known for some time. These devices are also known in the art as airbrasive or air-abrasive devices. Such devices vary in size and design depending on the particularly utility desired.\nOne area where use of these devices has proved advantageous is in the etching or abrading of small surfaces. Devices designed for this use are typically hand held and capable of delivering fine streams of air-abrasive material through narrow nozzles.\nA number of decades ago, the use of air-abrasive devices gained favor in the dental art The methods developed were termed \"airbrasive techniques\" and were designed to supplement the use of traditional dental drills to prepare a tooth for cavity repair, prophylaxis or other methods that required that a portion of the tooth be removed or that required the roughing of a tooth surface. The advantage of using air-abrasive techniques is that the dental patient experiences less trauma to the oral cavity due to the absence of perceptible pressure, vibration, noises created by the contact of a drill to tooth enamel, and heat created by frictional forces. This has resulted in reduced pain, apprehension, and fear by patients.\nOne disadvantage of the use of air-abrasive dental apparatus is that abrasive materials are dispersed into the oral cavity during use in a relatively uncontrolled fashion, can be inhaled by the patient, and are difficult to remove after a procedure is complete. Another disadvantage is that such particles can be dispersed into the air and create a hygiene problem. Abrasive particles can carry pathogens and blood particles from the mouth and permit those pathogens and blood particles to contact otherwise uncontaminated surfaces.\nSomewhat similar disadvantages exist with use of air-abrasive devices in other applications. Often it is desirable to prevent abrasive materials from contacting surfaces proximate to the target surface, from accumulating abrasive material on the target surface area, or from permitting fine abrasive particles from becoming airborne.\nSeveral devices have been developed to affect the dispersion of abrasive particles within the oral cavity. Coston, U.S. Pat. No. 5,197,876 discloses a splatter guard for air polishing dental devices. The guard comprises a bell-shaped flexible cone that is attached to the end of an air-abrasive device and guides abrasive particles towards the surface being treated. Ho, U.S. Pat. No. 5,356,292 discloses a dental sandblasting confiner in the form of a flexible transparent cup. The nozzle of a sandblasting device can be inserted in large opening of the cup which forms a mold around the nozzle. The Ho device contains additional openings for access to a tooth surface and for discharging output. Lokken, U.S. Pat. No. 4,611,992 discloses an anti-splash device that can be attached to a dental tool. The device comprises an inverted U-shaped member with legs for attaching the device to the dental tool. Wright, U.S. Pat. No. 4,850,868 discloses a spray shield comprising a modified tube that can be attached to the end of a dental handpiece. The device is used to direct material dispensed form the handpiece in a controlled fashion so as to minimize the amount of airborne particles.\nWhile the above cited inventions address one or more of the described disadvantages of air-abrasive systems, they are subject to several detrimental limitations. Although minimizing the amount of abrasive material released, by guiding it downward for instance, has certain benefits, it is more preferable to contain a substantial portion of released abrasive material and permit facile removal. Many of the devices in the prior art guide, but do not completely contain abrasive material nor permit easy removal thereof. Other devices that do permit removal of abrasive material are obtrusive and interfere with visualization of the surface to be abraded, making it difficult to perform precise dental procedures. Furthermore, those devices that do permit removal of abrasive material typically rely on a vacuum source to remove that material. Such a vacuum source adds additional expense and can also be intrusive.\nThus, there is a need for a device that can contain a substantial portion of the abrasive material expelled from an air-abrasive device while not obstructing visualization of the surface to be abraded and permitting removal of the abrasive material without the aid of a vacuum source."}
{"text": "Reverse osmosis is a well-known process for the separation of solvent from a solvent-solute solution. The process utilizes a semipermeable membrane and high pressures to selectively diffuse solvent molecules through the membrane resulting in separation of pure solvent from the solvent-solute solution.\nReverse osmosis systems typically utilize a cylindrical filter having walls formed from a semipermeable membrane material. A solvent-solute solution, from which pure solvent is to be removed, is contacted with the outer circumferential membrane wall of the filter. Upon application of sufficient pressure to the solution and filter exterior, i.e. typically at least about 100 psi and often from about 500 psi to about 1500 psi, solution is driven toward the interior of the filter. Pure solvent is then typically collected from the ends of the filter.\nFilter housings are provided for containing reverse osmosis filters. The housings are commonly formed from steel in view of the relatively high pressures existing within the housing during separation. As will be appreciated, steel housings are expensive and so have to some extent, limited the applicability and use of reverse osmosis systems.\nPlastic filter housings have been introduced which although satisfactory in many respects, still present several significant drawbacks. Currently available plastic filter housings for reverse osmosis filters utilize end caps that are bonded or otherwise affixed to the filter housing. Such attachment renders accessing the interior of the housing exceedingly difficult. Even if an end cap is successfully removed without damaging the cap or housing, such as to replace a filter cartridge, the end cap must then be re-attached. Re-attached end caps are often susceptible to failure upon pressurization of the filter housing. Thus, there is a need for an improved end cap assembly that enables ready removal of an end cap and secure re-attachment thereof to a plastic filter housing, without a loss in the integrity of the housing upon reassembly.\nAnother drawback to currently known filter housings, metal or plastic, is the limited degree of installation flexibility. That is, except for custom designed housings and systems which are typically relatively expensive, commercially available housings are available in only a limited number of connection configurations. Thus, it would be desirable to provide a filter housing with an increased number of connection configurations."}
{"text": "Implantable electrical stimulation systems have proven therapeutic in a variety of diseases and disorders. For example, spinal cord stimulation systems have been used as a therapeutic modality for the treatment of chronic pain syndromes. Peripheral nerve stimulation has been used to treat chronic pain syndrome and incontinence, with a number of other applications under investigation. Functional electrical stimulation systems have been applied to restore some functionality to paralyzed extremities in spinal cord injury patients. Stimulation of the brain, such as deep brain stimulation, can be used to treat a variety of diseases or disorders.\nStimulators have been developed to provide therapy for a variety of treatments. A stimulator can include a control module (with a pulse generator), at least one lead, and an array of stimulator electrodes on each lead. The stimulator electrodes are in contact with or near the nerves, muscles, or other tissue to be stimulated. The pulse generator in the control module generates electrical pulses that are delivered by the electrodes to body tissue."}
{"text": "The present invention relates to an assembly inspection system and method for examining the connected or bonded state of electronic parts of BGA (Ball Grid Array) type or CSP (Chip Scale/Size Package) type.\nIn recent years, in mounting an electronic part such as an LSI package on a substrate or board, a face bonding method is mainly adopted. BGA type electronic parts and CSP type electronic parts are typically face-bonding type electronic parts. Once an electronic part of the BGA type or that of the CSP type has been mounted in position, a soldered state thereof cannot be visually confirmed. In the process of assembling the board mounted with such a BGA type electronic part or a CSP type electronic part into an electronic equipment, there is a case in which it is impossible to decide whether an equipment failure is due to a defective connection of the BGA/CSP type electronic part or due to another cause. Thus, it is necessary to correctly inspect whether the BGA/CSP type electronic part is correctly bonded on the board. For this purpose, there is proposed an inspection system in which through-hole-provided pads are formed on a printed wiring board in correspondence with ball-shaped solder pads formed on a back surface of the BGA/CSP type electronic part, and in which the state of connection between the BGA/CSP type electronic part and the printed wiring board is inspected after the reflow process according to whether the through-hole-provided pads are filled with solder, as disclosed in Laid-Open Japanese Patent Publication No. 9-127178. Also, as disclosed in Laid-Open Japanese Patent Publication No.7-234189, there is proposed an inspection system utilizing a fine connection inspection device using a sensor.\nHowever, in the former of the above two conventional inspection systems, it is necessary to form on the printed wiring board the through-hole-provided pads corresponding to the ball-shaped solder pads on the back surface of the BGA/CSP type electronic part. Thus, the degree of freedom in designing a pattern deteriorates, which prevents a fine pattern design. In the latter of the two conventional inspection systems, it is necessary to introduce an expensive inspection device and thus the manufacturing cost becomes substantially high."}
{"text": "Sealed gas-tight containers require venting when gas pressure must be equalized between the interior and exterior of the container. Without venting, a flexible gas-tight container will bloat, leak, and possibly burst when the interior pressure exceeds the exterior pressure. Bloating can occur when the contents of the container generates gasses or heat by chemical reaction, for example when the contents include a peroxide-based toothpaste. Bloating can also occur when the container is stored in a heated environment. An unvented flexible gas-tight container will collapse when the internal pressure is reduced, for example when atmospheric oxygen is scavenged by one of the ingredients housed in the container. Some modes of transportation put a container at risk of both bursting and collapse. During vehicular transport through mountains and valleys, for example, a container is subjected to pressures that can rise above and drop below sea level pressure. Rigid gas-tight containers, such as glass containers, are susceptible to bursting or imploding if the internal and external pressures become sufficiently discrepant.\nOne way to equalize pressure is to provide a filter vent in the cap, lid, or other closure of a container. Filter vents generally include a gas-permeable filter, or other gas permeable microporous medium, which is interposed between the interior of the container and a vent aperture. The filter vent permits gases to diffuse in and out of the interior of the container, via the vent aperture, while excluding particulates larger than a threshold size, as well as liquids of a particular range of hydrophobicity. The term “filter cap” will be used as a generic term for all container closures fitted with a filter vent.\nFilter caps such have a serious drawback. The pores of the filter are susceptible to being filled and clogged by the contents of the container. Existing solutions to this problem are not satisfactory. PCT Application No. WO 97/02994 to Vakharia (“Vakharia”) discloses a gas permeable membrane mounted at an angle to the inner upper surface of a bottle cap. Angling of the membrane is intended to cause liquid to roll off the membrane surface instead of remaining fixed to the membrane surface. U.S. Pat. No. 6,196,409 to Lake et al. (“Lake”) discloses a cap or container including a filter vent that has no structural adaptation to prevent the filter from contacting the contents of the container. Instead, Lake relies on the matching of the filter material to the physicochemical properties of the contents of the container, so that the filter material encourages the phase separation and run-off of contents that have contacted the filter. This material-matching requirement limits the selection of filter materials that can be used with any particular type of contents. Furthermore the encouragement of drainage of container contents after they have contacted a filter is less desirable than the prevention of contact between the contents and the filter in the first place.\nContainer-mouth baffles, such as sifters, are used to regulate the dispensing of granular materials through the mouth of a container. A typical example is the sifter top disclosed by U.S. Pat. No. 5,513,781, to Ullrich, et al. (“Ullrich”). The sifter top includes a perforated liner extending across the mouth of a bottle suitable for containing spices or similar materials. Sifters typically include perforations large enough to permit the passage of granular materials out of a container.\nContainer-mouth baffles can potentially protect the filter of a filter cap from the contacting the contents of a container, but they have never been situated, or provided with appropriate perforations, to perform such a protective function. There is a need for assemblies and methods for protecting a filter cap with a container-mouth baffle."}
{"text": "Support matrices of the type given above are typically in form of porous monolithic beds or porous or non-porous particles. The support particles may be used in fluidised beds or packed beds. In order to reach a sufficient rigidity, the matrices often have to be cross-linked.\nTypical problems associated with support matrices for the above-mentioned uses are:                (a) high flow velocities may lead to an increased back pressure which ultimately causes collapse of the matrix;        (b) when in particle form, the matrices may have a tendency not to withstand mechanical forces during packing in columns and other reaction vessels;        (c) tricky to manufacture; and        (d) the matrices may need a separate sterilization.        \nOne previous way of manufacturing support matrices that have an improved rigidity is to start from a polymer carrying nucleophilic groups and groups that have a masked electrophilic reactivity activating the masked groups and permitting the electrophilic and nucleophilic groups to react with each other. The polymer is typically a polysaccharide. See WO 9738018 (Amersham Pharmacia Biotech AB).\nAnother recent method is to heat a polymeric support matrix carrying residual unsaturated groups. This method is in particular adapted to support matrices obtained by polymerisation of mixtures comprising vinyl aryl monomers, aryl or methacryl monomers or other unsaturated monomers. See International Patent Application PCT/SE99/0126 (Amersham Pharmacia Biotech AB).\nWO 9726071 (Amersham Pharmacia Biotech AB) and WO 9700255 (Amersham Pharmacia Biotech AB) include electron beam irradiation and gamma-ray irradiation as two out of many examples for initiating polymerisation of unsaturated groups in order to manufacture support matrices for various uses.\nEP 534016 (W.R. Grace & Co) describes chromatographic matrices with thermoresponsive properties. The publication suggests that one could cross-link beads consisting of linear polymers by irradiation with light, electron rays and gamma rays. It is also suggested that linear polymer after adsorption to a base matrix could be crosslinked in the same way.\nWO 9603147 (Fidia Advanced Biopolymers) suggests production of gels and biocompatible substrates by electron beam irradiation of a polymer funtionalized with a plurality of unsaturations. The uses are as carriers in drug administration, as biocompatible substrates for cell culture in three-dimensional systems etc.\nEP 159694 (E.I. Du Pont de Nemours) suggests electron beam irradiation for the production of gels from monomers with one and two polymerisable unsaturations.\nThere are no demands on withstanding liquid flow on the monolithic or packed bed formats of the matrices described in WO 9603147 and EP 159694."}
{"text": "1. Field of the Invention\nThe invention relates to an aerofoil having a slat mounted to extend and retract relative to the aerofoil to increase lift especially during landing and take off of an aircraft.\n2. Discussion of Prior Art\nIn certain aircraft, it is usual to provide a data link between the slat and the aerofoil to enable data such as pressure and temperature to be monitored. Known data links necessitate use of the usual slat mounting track to form a component of the data link. That involves complexity and considerable expense and an object of the present invention is to provide an aerofoil having a slat with an improved data link arrangement."}
{"text": "Work vehicles, such as tractors and other agricultural vehicles, typically include one or more electronic control units (ECUs) configured to receive input signals and/or transmit output signals for monitoring and/or controlling vehicle operation. For example, work vehicles often include an emission sensor control unit communicatively coupled to one or more emission sensors, such as a nitrous oxide (NOx) sensor, configured to measure the concentration of one or more emission gases contained within the exhaust flow expelled from the vehicle's engine. As such, the emission sensor control unit may be configured to receive emission-related measurement signals from the sensor(s). Suitable control signals may then be transmitted from the emission sensor control unit to an engine control unit of the work vehicle to allow the operation of the engine to be controlled based on the emission-related measurements provided by the sensor(s).\nIn general, it is often desirable to position the control unit(s) of a work vehicle under the vehicle's hood to allow such control unit(s) to be in close in proximity to the components with which the unit(s) must communicate. However, electronic control units are typically required to be maintained below a specific operating temperature. Unfortunately, at high ambient temperatures and/or during high load applications, the under-hood air temperature of a work vehicle may often reach and/or exceed the predetermined temperature limit defined for a given control unit. As a result, there is an increased likelihood of damage occurring to such control unit due to overheating. In addition, after a hot shutdown, the under-hood air temperatures of a work vehicle often increases well above the under-hood temperatures typically experienced during vehicle operation, thereby further increasing the likelihood of damage occurring to any control units installed under the hood.\nMoreover, conventional under-hood mounting arrangements typically require that a control unit be mounted to a separate mounting bracket installed under the hood. Such mounting brackets unnecessarily increase the number of under-hood components required for a work vehicle. In addition, since control units typically need to be isolated from vehicle vibrations to operate properly, vibration isolators must often be associated with the mounting bracket(s) to satisfy the vibration requirements of such control units, thereby increasing the overall costs of the installation.\nAccordingly, an improved under-hood mounting configuration for a control unit of a work vehicle that addresses one or more of the issues highlighted above would be welcomed in the technology."}
{"text": "Various antiphlogistic and analgesic agents are known but it is generally found that the antiphlogistic and analgesic agents conventionally used have a rather high toxicity and often induce various undesired side effects such as for example gastric lesions. There is thus a need for antiphlogistic and analgesic agents having a lower toxicity and no undesired side effects.\nIt is known that chloro-5-benzo(b)furyl-2 carboxylic acid, chloro-5-benzo(b)furyl-2-acetic acid, methoxy-5 -benzo(b) furyl-2-acetic acid and a derivative of substituted .beta.-(benzo(b)furyl-2) acrylic acid have a lower anti-inflammatory activity and a minor toxicity than phenylbutazone (Eur. J. Med. Chem. - Chimica Therapeutica, March-April 1975-10, N. 2, p. 182-186).\nThe present invention is concerned with the following 5-substituted-3-methyl-2-benzo(b) furylacetic acids of general formula ##STR1## wherein R represents a hydrogen atom; an alkyl group containing 1 to 6 carbon atoms, e.g. a methyl, ethyl, n-propyl, isopropyl, n-butyl, sec-butyl, tert-butyl, n-pentyl or n-hexyl group; an alkoxy group containing 1 to 6 carbon atoms, e.g. a methoxy, ethoxy, n-propoxy or isopropoxy group; a cycloalkyl group containing 3 to 6 carbon atoms, e.g. a cyclopropyl, cyclobutyl, cyclopentyl or cyclohexyl group; a cycloalkoxy group containing 3 to 6 carbon atoms, e.g. a cyclopropoxy group; an alkenyl group containing 3 or 4 carbon atoms, e.g. a propenyl or 1-butenyl group; an alkenyloxy group containing 3 or 4 carbon atoms, e.g. a propenyloxy or isopropenyloxy group; a cyclohexenyloxy group; a phenyl group or a substituted phenyl group having the formula ##STR2## wherein R.sub.1 represents a methoxy group; a halogen atom, e.g. a fluorine or chlorine atom; or an alkyl group containing 1 to 4 carbon atoms, e.g. a methyl, ethyl, propyl, n-butyl, sec-butyl or tert-butyl group; a phenoxy group; a benzyl group; a substituted benzyl group having the formula: ##STR3## wherein R.sub.1 is as defined above; a trifluoromethyl group; or a halogen atom, e.g. a fluorine, chlorine or bromine atom and phsyiologically acceptable salts thereof.\nThe present invention is further concerned with the discovery that the 5-substituted-3-methyl-2-benzo(b)furylacetic acids of general formula (I) show antiphlogistic and analgesic activities.\nThus according to the present invention there are provided pharmaceutical compositions comprising as active ingredient a compound of formula (I) as hereinbefore defined or a physiologically acceptable salt thereof in association with a pharmaceutical carrier or excipient.\nAmong the derivatives of benzo (b) furylacetic acid represented by the formula (I), the derivatives where R is a hydrogen atom [Chem. Ber. 97, 3577 (1964)], a methyl radical [Rocz. Chem. 44, 1913 (1970)] and a chlorine atom [Indian Journal of Chem., 2, 456 (1964)] were previously known, although it was not known that these compounds had antiphlogistic and analgesic properties. The other compounds of general formula (I) are however new compounds and constitute a further feature of the present invention.\nThe preparations of the three known compounds, 3-methyl-2-benzo(b) furylacetic acid, 3,5-dimethyl-2-benzo(b)furylacetic acid and 5-chloro-3-methyl-2-benzo(b)furylacetic acid are described in the above indicated literature references. However, these preparations have disadvantages of complicated steps and a low yield.\nThe present invention is accordingly to provide processes for preparing the derivatives of benzo(b)furylacetic acid of the formula (I) with a high yield in a simple manner. According to the present invention, the derivatives of benzo(b)furylacetic acid represented by the formula (I) can be prepared by the process exemplified as follows:\n(A) Process using as a starting material 3-methyl-2-cyanomethylbenzo(b) furan or 5-substituted derivatives thereof having the general formula (II) ##STR4## wherein R is as hereinbefore defined.\nThe compounds of general formula (I) can be prepared from the starting compound by hydrolysis in an aqueous medium in the presence of an acid or alkali.\nAcids which are preferably used for this purpose are, for example, hydrochloric acid, sulfuric acid, phosphoric acid and para-toluenesulfonic acid. It is also possible to use a strongly acidic ion exchange resin.\nAlkalies which are preferably used for this purpose are, for example, hydroxides of alkali metals and alkaline earth metals, e.g. sodium hydroxide, potassium hydroxide and calcium hydroxide.\nThe hydrolysis is preferably effected in the presence of a solvent, especially an organic solvent which does not react with the starting compound and which is miscible with water. Preferred organic solvents include, for example, dioxane, tetrahydrofuran, acetone, 2-methoxyethanol, 2-ethoxyethanol and lower alcohols, e.g. methanol, ethanol, n-propanol, isopropanol, n-butanol and tert-butanol. Advantageously the organic solvent is present in an amount of from 1/4 to 4 times the volume of water. In the reaction, 0.01 mole to 1.0 mole, preferably 0.1 mole-0.5 mole, of the starting compound is used per one liter of the total solvent.\nThe reaction is usually carried out at a temperature within the range of from 50.degree. to 100.degree. C., preferably at about the boiling point of the organic solvent used. After completion of the hydrolysis, the reaction mixture is concentrated. If desired, the reaction mixture is treated with active charcoal prior to the concentration. The pH of the concentrated mixture is then adjusted to a pH of 2 with an acid, for example, hydrochloric acid or sulfuric acid, and is cooled to precipitate crude crystals. The crude crystals are then recrystalized from a suitable solvent, for example an aqueous solution of an alcohol to yield the purified final compound.\nThe starting compounds of general formula (II), 3-methyl-2-cyanomethylbenzo(b)furan and 5-substituted derivatives thereof, can be prepared, for example, by the process disclosed in DT-AS 1,203,277. They can also be prepared as follows:\nPhenol or an appropriate p-substituted-phenol of formula ##STR5## wherein R is as hereinbefore defined, is reacted with an equimolar amount of sodium hydroxide in an aqueous medium to give a sodium salt of phenol or the p-substituted phenol respectively which is then treated with ethyl .alpha.-chloracetoacetate to yield a compound of formula ##STR6## wherein R is as hereinbefore defined. The ethyl .alpha.-(phenoxy)acetoacetate or ethyl-.alpha.-(p-substituted-phenoxy) acetoacetate obtained is then subjected to ring-closure by treatment with concentrated sulfuric acid to give 3-methyl-2-ethoxycarbonylbenzo(b) furan or a 5-substituted derivative thereof of formula ##STR7## wherein R is as hereinbefore defined. The resultant compound is subjected to reduction with lithium aluminium hydride in an organic solvent, e.g. diethyl ether, to give 3-methyl-2-hydroxymethylbenzo(b)furan or a 5-substituted derivative thereof of formula ##STR8## wherein R is as hereinbefore defined. This compound is then treated with thionyl chloride to give 3-methyl-2-chloromethylbenzo(b) furan or a 5-substituted derivative thereof of formula ##STR9## wherein R is as hereinbefore defined, which is then subjected to reaction with sodium cyanate to give 3-methyl-2-cyanomethylbenzo(b) furan or a 5-substituted derivative thereof of formula (II).\n(B) Process using as a starting compound 3-methyl-2-acetylbenzo(b)furan or 5-substituted derivatives thereof having general formula (III) ##STR10## wherein R is as hereinbefore defined.\nThe starting compound, which can be obtained, for example, by the process described in Bull. Soc. Chim. Fr., 1970 (10), page 3601, is reacted with hydrogen sulfide or sulphur together with ammonia, a primary amine or a secondary amine in an aqueous solution or in an aqueous solution of organic solvent to obtain its thioamide.\nAmines which may be used for this purpose include, for example, methylamine, pyridine, morpholin and the like.\nThe organic solvent which may be used for the reaction includes, for example, dioxane, piperidine and the like. It is also possible to use the amine as both the reaction material and the organic solvent. In the reaction, 0.01 mole to 2.0 mole, preferably 0.1 mole-1.5 moles of the starting compound is used per one liter of the total solvent. When morpholin is used as both the reaction material and the organic solvent, as for the amount of sulphur and morpholin, 1.3 to 1.7 moles, preferably 1.5 moles of both compounds are used per one mole of the starting material of general formula (III).\nThe reaction is preferably carried out at a temperature within the range of from 130.degree. to 200.degree. C. for about from 5 to 25 hours. After completion of the reaction, the reaction mixture is cooled to precipitate crystals of the thioamide. If desired, it is also possible to otain the thioamide with a better yield by concentrating the reaction mixture, removing the solvent and adding water.\nThe thus-obtained thioamide is subjected to the hydrolysis with a suitable acid or alkali to give the desired 5-substituted-3-methyl-2-benzo(b)furylacetic acid represented by general formula (I). The acids which may with advantage be used for this purpose are exemplified by hydrochloric acid, sulfuric acid, paratoluene-sulfonic acid and the like; and the alkalis which are preferably used for this purpose include for example hydroxides of alkali metals and alkaline earth metals such as sodium hydroxide, potassium hydroxide and calcium hydroxide. In the reaction, 0.01 mole to 1.0 mole, preferably 0.1 mole-0.5 mole of the starting compound is used per one liter of the total solvent.\nThe desired final product can be isolated for example in the following manner.\nThe reaction mixture containing the desired final product obtained by the hydrolysis is adjusted to a neutral pH for example with sodium carbonate and then washed with a suitable solvent such as for example benzene. After this, the water layer in the reaction solution is adjusted to an acidic pH 2, for example with a suitable dilute acid such as dilute hydrochloric acid, dilute sulfuric acid and the like to isolate the desired acid which is then extracted for example with ethylacetate and the like. After removal of water from the extract, the organic solvent of the extract is concentrated to obtain the desired crude product. The recrystalization is then carried out by the use of a suitable solvent such as, for example, an aqueous solution of alcohol, ethanol and the like.\nPreferred physiologically acceptable salts of 3-methyl-2-benzo(b)furylacetic acid and the 5-substituted derivatives thereof of general formula (I) for incorporation into the pharmaceutical compositions according to the invention are, for example, the sodium, potassium, calcium and aluminium salts (obtained by reaction with the appropriate metal oxide or hydroxide) and also the physiologically acceptable non-toxic ammonium and amine salts (obtained by reaction with for example ammonia, diethylamine or triethanolamine).\nThe pharmaceutical compositions according to the present invention may with advantage be presented in a form suitable for internal or external administration, e.g. oral or topical administration, for example as tablets, granules, powders, capsules, ointments or creams which may be prepared in conventional manner.\nSuitable daily dosages of the active compounds of general formula (I) whether administered internally or externally are generally in the range from 50 to 600 mg per man per day, depending upon the symptoms of the patient.\nThe compositions are preferably presented in the form of dosage units, each dosage unit being adapted to provide a fixed dose of active ingredient. Such compositions preferably contain from 25 to 100 mg of active ingredient per dosage unit.\nCompositions in the form of tablets may be formulated with suitable amounts of various other ingredients such as, for example, excipients (e.g. lactose, glucose, sucrose or mannitol), disintegration agents (e.g. starch, sodium alginate, calcium carboxy-methyl-cellulose (CMC), crystal cellulose or sugar ester), lubricants (e.g. magnesium stearate or talc) and binding agents (e.g. simple syrup, gelatine solution, polyvinyl alcohol or polyvinylpyrrolidone). All of the abovementioned ingredients are used in the amounts conventionally used in the preparation of tablets. Dispersants (e.g. methylcellulose) and plasticizers may be used as coating agents for tablets.\nCompositions according to the present invention presented in forms other than tablets and capsules, e.g. powders, preferably contain the active compounds of general formula (I) in amounts of from 5-10% by weight of the composition, together with other ingredients, generally excipients, e.g. lactose, glucose, sucrose and mannitol.\nThe acute toxicity, antiphlogistic and analgesic activities and the effect on gastric mucosa of 3-methyl-2-benzo(b) furylacetic acid and certain of the 5-substituted derivatives of general formula (I) according to the invention have been examined as described hereinafter. In the tests, the specimens Nos. 1-8 are the active compounds represented by the formula (I) according to the present invention and the specimens Nos. 9 and 10 are known compounds having known antiphlogistic and analgesic activities included for comparison purposes."}
{"text": "In today's modern society privacy issues have become a major concern. Modern technologies such as computers and the Internet enable people to access more information at rates never before possible. However, many of the tasks for which people use the Internet are considered highly private or confidential matters. For example, the Internet has streamlined tasks such as banking, filing tax returns, conferencing with business associates, and communicating with friends and family.\nCoupled with the advantages that accrue from the Internet is an increased susceptibility to malicious eavesdropping and/or cyber-attack. Thus, as the tools with which people conduct their daily affairs advance in complexity, so too must the means by which private or confidential matters are concealed and the communication tools safeguarded from evildoers and/or failures. As such, industry leaders have organized a Trusted Computing Group (“TCG”) to address these concerns.\nTCG is a an industry standards body, including computer manufacturers, device manufacturers, and software vendors who have a stake in promoting the security of computing platforms and devices (see, https://www.trustedcomputinggroup.com). One goal of TCG is to promote a security chip called the trusted platform module (“TPM”). The TPM is defined by specifications published by the TCG, such as the TCG Main Specification, Version 1.1b, Feb. 22, 2002. The TPM is an isolated silicon device built into the motherboard of a personal computer (“PC”), but separate from the central processing unit (“CPU”), for establishing trust and trust metrics in a Trusted Computing Environment. “The TCG architecture and the TPM enable platform level authentication, attestation, and the generation and storage of cryptographic keys. The TPM can also perform certain operations on these keys such as signing and hashing within its protected locations. The TPM specifications are designed to enable this increased level of trust and security while maintaining privacy, user choice and control.” TCG, Corp., An Overview Of The TPM And The TCG's Trusted Platform Architecture, Version 0.99D, (May 5, 2003).\nSoftware vendors, such as Microsoft, have begun developing software platforms that will take advantage of the security features provided by a TPM silicon device. For example, a future version of Windows called Longhorn is expected to provide numerous security functions, including encryption, privacy safeguards, and digital rights management with use of the TPM silicon device.\nHowever, widespread acceptance and implementation of the TCG architecture and the TPM is still uncertain. Furthermore, integration of the TCG architecture with network servers, where it is crucial to provide reliable, available, and serviceable functionality, has proven to be difficult. For example, some of the functionality defined by TCG architecture has been found to conflict with current techniques for “hot plugging” network servers."}
{"text": "Modern motor vehicles commonly feature occupant safety systems including, for example, seat belts, frontal collision airbags, side airbags, and curtain airbags. One or more sensors on-board the vehicle feed signals to an electronic control module which, based on the signals, detects a collision or other unsafe condition so that restraints and/or other safety systems may be activated.\nA typical passenger vehicle has two rows of seats for the driver and passengers, and such vehicles often are equipped with a side impact sensor located in, on, or adjacent to the door or sidewall of the vehicle on either side of each seating row, for a total of four sensors (two on each side of the vehicle). The side impact sensors are typically accelerometers or pressure sensors. Such a sensor suite is effective for sensing lateral collisions or impacts that occur on the sides of the seating rows.\nExtending this concept to larger vehicles that may have three or more passenger seating rows would require six or more total side impact sensors."}
{"text": "a. Field of the Invention\nThe present invention generally relates to a mobile communication device capable of a wireless communication, such as a so-called mobile telephone (or cellular phone), a personal handy phone system (PHS), etc., and particularly to such devices having a flat display of a relatively large size.\nb. Related Art\nIn recent years, mobile communication devices such as a mobile telephone and a PHS are widely spread while the number of users of mobile PCs and personal digital assistants (PDAs) are rapidly increasing. In this connection, it becomes possible to perform a data communication outdoors or at an outside site with any desired person by connecting the mobile PC with the mobile communication device.\nA product has recently been commercialized in which the two devices are integrated, i.e., a mobile communication device in the form of a PHS telephone with a PDA capability. This product could also be considered as a mobile information device in which a PDA is provided with a PHS capability.\nThis type of mobile communication devices are available for a person at a place other than his or her home or office not only for simply as a telephone but also for transmit ting/receiving e-mails and facsimiles, or accessing Internet WWW (world Wide Web) servers to browse home-pages.\nFor this end, such mobile communication devices include a display of a larger screen size than that of conventional mobile telephones.\nIn making a phone call, the dialing can be performed not only by pressing numeral keys, but also based on the phone number of a person selected from a telephone directory data, which is stored in the device, without the key entering of the phone number.\nOn the other hand, information of home-pages obtained from the WWW servers is written in Hyper Text Markup Language (HTML). A document written in HTML, itself, is a text file in which codes sandwiched by symbols \"&lt;\" and \"&gt;\", called tags, are used to specify character-attribute information or layout information. A program called a \"browser\" interprets the tags so as to display the HTML document.\nAn HTML document can have other text, image or audio inserted therein, and also have parts thereof linked with images or other HTML documents. That is, in HTML documents, a location (i.e., an anchor point or hot-spot explained below) in a document can be linked (or correlated) with another location or a separate file, etc. This enables an immediate movement from a location to another location of a linked destination, allowing the effective browsing of a variety of documents (including so-called multi-media). In addition, it is possible, based on a selection of an anchor point, to open a window for sending an electronic mail (e-mail) or to activate a particular function that is associated with the anchor point.\nThus, anchor points set in an HTML document with respect to various items therein allow a user to select them so as to view other documents or activate various functions which are associated with the selected anchor points.\nIncidentally, most of the portable information devices such as a PDA are provided with pens, which are used as a pointing device to directly enter commands through a touch panel on a display of the device.\nHowever, when using the mobile communication device simply as a telephone, it is bothersome to use the pen each time a user specifies one of the phone numbers in a telephone directory. Usually, therefore, motion keys are used for selectively indicating two directions to search a phone number. For example, phone numbers with anchor points set in a telephone directory written in HTML can sequentially be designated by the use of the motion keys or the like. The browser can recognize the position of the anchor points so that for selecting one anchor point based on the motion key operation, the anchor points are sequentially designated one by one simply by indicating a forward direction or a backward direction with the keys. Such designation of an anchor point is referred to as a \"focus\" in the specification. Also, the position of an anchor point that is currently focussed in a displayed HTML document is referred to as a \"focussed\" position. An operation to cause an execution of a particular function or a jump associated with such focussed anchor point is referred to as a \"selection. The \"selection\" is distinguished from the \"designation\" which is a step taken prior to the selection.\nIn the above-mentioned mobile communication device with a PDA capability, it is preferable to include in the telephone directory not only names and telephone numbers but facsimile numbers or e-mail addresses, etc. Further, with the telephone directory being an HTML document, the additional items can be provided with various functions as described above.\nHowever, in such a case where anchor points are set to various items in the HTML document, the designation of an anchor point by the use of the motion keys mentioned above suffers from a problem in which it is difficult to quickly find a desired telephone number because a lot of undesired anchor points could be encountered until an anchor point of the desired number is reached.\nThe present invention was made in view of the foregoing problem, and it is an object of the invention to enhance the operability of a mobile communication device with a flat display of a relatively large screen size, capable of displaying thereon an HTML document.\nIt is another object of the invention to provide a mobile communication device with a rotary push switch, in which a desired anchor point for automatic dialing is quickly reached in response to an operation for moving a focus."}
{"text": "1. Field Of The Invention\nThe present invention relates generally to back pressure valves, and more particularly, but not by way of limitation, to back pressure valves suitable for use in a workover string for acidizing of wells.\n2. Description Of The Prior Art\nDuring certain workover jobs on oil and gas wells, a tool known as a back pressure valve is run in the work string. The back pressure valve typically is a flapper type valve which when in a closed position prevents upward flow through the work string. The flapper is opened by a stinger tube which is pushed downward through the flapper. The operation of the stinger tube is typically controlled by a J-slot interconnection between the stinger tube and the housing which contains the flapper valve. Thus, the flapper valve is opened by the weight of the workover string above the back pressure valve which is available to be set down on the stinger tube to force the flapper open. Sometimes a difficulty is encountered when there are very high pressures contained in the workover string below the flapper. The back pressure valve is typically placed at a relatively high point in the work string near the surface, and thus the weight of the work string thereabove is sometimes inadequate to force the flapper open against the high internal pressures trapped below the flapper.\nThus, there is a need for a back pressure valve which can be operated reliably regardless of the pressures acting against the flapper and holding it closed."}
{"text": "Kaolin is a naturally occurring, relatively fine, white clay which may be generally described as a hydrated aluminum silicate. Kaolin clay is widely used as a filler and pigment in various materials, such as rubber and resins, and in various coatings, such as paints and coatings for paper.\nSlurries of kaolin clay are generally made by mixing a quantity of kaolin clay with water and other optional additives. However, chemicals added during beneficiation to achieve various desired results in these slurries, such as dispersants, colloidal thickeners, ammonia and leaching chemicals, also may provide usable energy sources for aerobic and anaerobic microbial growth. Because of these additives, most commercial kaolin slurries are susceptible to microbial spoilage and, therefore, require high levels of microbiocides to control aerobic and anaerobic microbial growth. Examples of such microorganisms include algae, bacteria, mold, spores, etc. The growth of microorganisms and their by-products tends to adversely affect the properties of a kaolin clay slurry by altering the color, odor and viscosity of the slurry. In many instances, the altered slurry may have little or no commercial value.\nThe problem of such microbial contamination in kaolin clay slurries is similar to problems in other industries, such as the pulp and paper industry or the petroleum industry, which have aqueous sources that encourage microbial growth. Microbial contamination can also significantly affect the efficiency of industrial structures or processes such as cooling towers or lubricating systems. Because of these problems, microbiocidal agents (i.e., microbiocides) are conventionally used, often at very high dosages, to eliminate or reduce the growth of microorganisms.\nMicrobiocides are well-known agents for the control of microbial growth in aqueous systems, such as kaolin clay slurries. However, the use of such agents raises environmental concerns which influence the choice of a microbiocide. Many of the early microbiocides were chlorinated or mercurial compounds which left harmful by-products in the environment. Some of the formerly used compounds, such as formaldehyde, are no longer environmentally acceptable because these compounds were found to be carcinogenic or teratogenic. These problems, therefore, limit the number of microbiocides available for use to those microbiocides which are less effective and more expensive.\nThere are a variety of problems associated with the use of microbiocides in the kaolin clay industry. These problems include the environmental impact associated with many microbiocides, the high cost of microbiocides and the large amounts of microbiocides that may be necessary to reach the desired microbial control level in a specific situation.\nMicrobiocides which are commonly used in kaolin clay slurries include 1,5-pentanediol; tetrahydro-3-5-dimethyl-2H-1,3,5-thiadiazine-2-thione; 1,2-benzisothiazolin-3-one; and 5-chloro-2-methyl-4-isothiazolin-3-one/2-methyl-4-isothiazolin-3-one. The disadvantages to the use of these microbiocides include the high cost of their use at the necessary higher dosages, their ineffectiveness in certain slurries and their environmental impact.\nAs previously indicated, slurries of kaolin clay differ widely in their susceptibility to microbial contamination. Process additives are key factors but other factors, such as the degree of mechanical or chemical processing, variations in crude clay sources, differences in crude clay particle size distributions and crude clay impurities, may also create additional problems.\nA method to solve some of the problems associated with microbial contamination is disclosed in Friedman et al. U.S. Pat. No. 4,975,109, which provides for using a combination of chemicals to act in a synergistic fashion to kill microorganisms. However, this patent requires the addition of an oxidizing agent, a microbiocide, a surfactant and an anti-corrosive material.\nFor various reasons, the prior art fails to provide an acceptable process that will control the level of microbial activity in a variety of kaolin clay slurries without the cost related to high dosages of the combinations of chemicals. Therefore, a need exists in the kaolin clay industry for an effective and broadly applicable process to reduce microbial growth and, therefore, provide commercially useful kaolin clay slurries having greater utility (i.e., a longer shelf life)."}
{"text": "The invention is based on a method and an apparatus as generally set forth hereinafter. Means for controlling the spring stiffness, for instance of a motor vehicle, are known (German patent No. 16 30 058); in this known apparatus, two work chambers of a shock absorber or telescoping spring are connected via external lines to an apparatus comprising a pump and two reservoirs. Only one-way check valves are disposed in the connecting lines to the telescoping spring. However, with this type of apparatus, the damper stiffness of this kind of shock absorber cannot be varied, because to do so energy must be supplied from outside--via the pump--which takes a relatively long time and means that there is a certain energy demand. Controlling the damper stiffness in a shock absorber is also known from German Offenglegungsschrift No. 33 04 815. Suspension systems of present-day vehicle types, in particular passenger vehicles, are typically optimized to an average operational case in terms of the spring stiffness and damper stiffness, with parameters being structurally fixed and remaining unchanged, except for effects associated with aging, during driving. Since in extreme operational cases, such as an empty or fully loaded vehicle, or with varying vehicle movement parameters (rapid cornering, braking, acceleration, smooth highway driving, and the like), optimal suspension or damping of the suspension system is not attained in all such operating situations, it is also already known to switch over among a plurality of damper settings or spring stiffnesses. This prevents long-term adaptation of the suspension system (and possibly the damper system, which is either arbitrarily integrated into the suspension system itself or is a component thereof), and especially it is impossible to make an automatic, finely-tuned adjustment to various road conditions or various kinds of driving, because a switchover in the characteristic curves of the suspension and/or damping system can be made only in stages, typically between only two operating states.\nIt is also known (U.S. Pat. No. 3,807,678), in a suspension system involving two masses, one of which may be one or more wheels of a vehicle and the other the vehicle body, to dispose a standard, passive compression spring between the two masses, which is called a passive isolating element and has a so-called active damper switched parallel to it. This damper, in which a piston slides in a cylinder and divides it into two work chambers, is considered to be active because an intervention is made into the damping properties, that is, into the positive volume displacements of the pressure medium in the various working halves of the damper by control means, in a so-called active manner. To this end, the two working chambers are connected crosswise and parallel to each other via opposed valves allowing a flow of pressure medium in only one direction; the amount of pressure medium then allowed to pass through these valves then also becomes \"active\" by appropriate control of the valves by means of suitably prepared sensor signals. Because in this known suspension system, the spring itself is entirely passive, but the damper is conceived of as being active in terms of its properties, the overall system in this patent is called a semiactive system. However, this term is not semantically related to the dampers of the present invention, which without reference to suspension systems not taken into account are themselves designated as so-called semiactive dampers, for reasons to be explained hereinafter.\nIt is also known, in wheel suspensions in vehicles, to provide so-called active damping means (see the article, \"Active Damping in Road Vehicle Suspension Systems\", published in the periodical, Vehicle System Dynamics, 12 (1983), pages 291-316). This publication is referred to also because it includes basic concepts, in theoretically detailed form, applicable in particular to active damping properties."}
{"text": "The present invention relates to a microfabrication technique which may be advantageously used, for example, for the manufacture of magnetic bubble memory devices, semiconductor devices, and the like. More particularly, the present invention relates to a method for forming a pattern of metal elements arranged with very small gaps therebetween on a substrate of a device.\nIn well-known field access types of magnetic bubble memory devices, a permalloy pattern for the propagation of magnetic bubbles (permalloy propagation pattern) is formed on the bubble-supporting magnetic layer of a substrate. The permalloy propagation pattern is constituted by a plurality of pattern elements made of thin films of permalloy. These permalloy pattern elements are chevron-shaped or half disc-shaped, for example, and are arranged with very small gaps therebetween. In a magnetic bubble memory device using 2 .mu.m diameter bubbles, the size of the gap is approximately 1 .mu.m.\nThe pattern of metal elements such as the permalloy propagation pattern described above is conventionally formed by using photolithographic techniques. That is, first, a pattern forming layer of metal is formed on a substrate of a device. Next, the pattern forming layer is coated with a layer of photosensitive resist material (photoresist). An optical exposure process is used to form a resist pattern corresponding to the pattern of metal elements to be formed from the layer of photoresist. Thereafter, an etching process is used to form the pattern of metal elements from the pattern forming layer.\nHowever, the smallest size of the resist pattern which can be formed on the pattern forming layer by using the photolithographic technique is approximately 1 .mu.m. Therefore, according to the conventional method, it is impossible to form a pattern of metal elements with gaps of a size smaller than 1 .mu.m.\nIn recent years, there have been a number of proposals for realization of a high density magnetic bubble memory device using 1 .mu.m diameter bubbles. Such a device requires a permalloy propagation pattern with gaps of a size less than 0.5 .mu.m. In order to satisfy this requirement, the electron-beam lithographic technique or X-ray lithographic technique is used. However, electron-beam systems and X-ray systems are large in size and high priced, thereby increasing the overall cost of the equipment."}
{"text": "In general terms, the present invention relates to a self-propelled drawing device or a tracing device which can be for example an output peripheral for a desk top computer. The drawing device can be referred to as an \"X-Y plotter\". During the drawing, it should be able to access substantially all of the drawing medium. The drawing device draws a desired trajectory or locus or suitable approximation thereto as defined by computing or calculating means. The data to be plotted is held as a series of vectors in computer memory. Each vector is converted to the necessary movements required to reproduce the data on the drawing medium. The conversion calculations from vectors to movements may be carried out in the desk top computer, or in the drawing device or in a secondary device which interfaces between the two. The tracing device works in reverse in that it is made to follow a line by manual propulsion, and signal the trajectory or locus of the line or a suitable approximation thereto to computing or storage means. Both the drawing device and the tracing device act as a drafting peripheral in the form of a locus-tracking device in that the drawing device tracks a locus in order to draw a line and the tracing devices tracks a locus on a planar medium."}
{"text": "The present invention relates to systems and methods for improving target aiming accuracy and/or target acquisition speed of participants during interactive video games.\nInteractive action video gaming systems have evolved and improved rapidly over the past twenty years. Many interactive video games on the Gen4 gaming systems, such as the Sony PlayStation4, are capable of providing players with an almost real-time interactivity. Many of these video games are action-based games involving the acquisition and shooting of target(s) while each player's avatar is on the move, and most of these video games present an aiming point proximate to a physical center of the gaming system's display screen.\nIn many of these video games, the aiming point disappears or changes shape when an avatar is on the move, and then reappears or reverts shape when the avatar is stationary. The absence or mutation of the aiming point during the avatar's movements substantially increases target acquisition and/or reacquisition time, and potentially also reduces the accuracy of the avatar. Further, crosshairs provided by the video games are often dim and hard to visualize.\nPrior feeble attempts at providing a more permanent aiming point on the display device include attaching a sticker or attaching an LED using a suction cup at the center of the screen. Both of these methods are deficient. A sticker does not work well in “darker” scenes, such as inside dimly lighted rooms. Disadvantages of attaching an LED with a suction cup include obscuring the area covered by the cup and also the area covered by the power supply leads for the LED. The physical size of the LED is another disadvantage.\nIt is therefore apparent that an urgent need exists for improved aiming capability for players of video games. This improved aiming capability will enable video gamers to quickly and accurately accomplish target acquisition while on the move."}
{"text": "Hair transplant procedures have been carried out for decades. Initially, a punch was used to remove a circular area of hairy skin containing ten or more follicular units (of 1-4 hairs each). The area of hairy skin replaced a like area of bald skin removed from the patient. Several of such “plugs,” were placed into areas in the bald part of the head.\nThe circular punch tool was later replaced by a hollow powered drill and the space left in the donor area was left to heal naturally. Both of these prior art procedures allowed wounds to stay open for weeks at a time exposing a patient to the discomfort from large wounds measuring 3-5 millimeters in diameter.\nToday there are two standard procedures for harvesting hair, the first involves a linear incision which permits the removal of a strip of hairy skin down into the fatty level of one quarter inch and measuring a number of square inches. The resulting wound is sutured closed and the strip is dissected into grafts (under a microscope), cooled in an ice bath or refrigerator and then transplanted into a bald area in needle size holes. Forceps grasp each graft and places them into holes in the bald area. In one form for dissection the hair from the strip of scalp uses blind harvesting of grafts from the strip of hairy skin which can result in significant damage to the hair. The damage occurs because the strip of hair in the hairy skin is forced through a cutting grid in order to make grafts of a predetermined size. The cutting blades of the grid are positioned at the most ideal distance between follicles. Unfortunately, the distance between follicles varies randomly. The result is that a significant number of the hair follicles can be damaged and die.\nThe second harvesting technique involves the use of a microscope to dissect the hair from the excised scalp.\nA third harvesting technique uses a punch which cores out from the scalp, the basic anatomical unit of hair, the follicular unit, which contains between one and four hairs each. The problems with this technique is that there can be considerable damage to the follicular unit as they are cored out from the scalp, one at a time because the skills to accomplish this are difficult to learn and master and the use of a sharp or dull punch to accomplish this runs into problems produced by the body's collagen which take on different consistencies with different people."}
{"text": "In many cases it is advantageous that equipment is modular so that the equipment comprises a body device and plug-in units installed in the body device. The equipment can be, for example but not necessarily, telecommunication equipment such as an internet protocol “IP” router, an Ethernet switch, an Asynchronous Transfer Mode “ATM” switch, and/or a MultiProtocol Label Switching “MPLS” switch. A commonly used construction is such that the telecommunication equipment comprises a frame and plug-in units which are installed in plug-in unit slots of the frame. In this case, the frame represents a body device for receiving the plug-in units. Electrical connectors in a plug-in unit make galvanic contacts with corresponding electrical connectors in the frame when the plug-in unit is inserted in the plug-in unit slot of the frame. The frame may have wirings such that plug-in units installed in the frame form a full mesh network or such that plug-in units installed in the frame are connected to a central element which may comprise one or more plug-in units installed in the frame or which may be a functional component integrally built in the frame.\nCare is needed in the installation of a plug-in unit in a body device because galvanic contacts between the electrical connectors of the plug-in unit and corresponding electrical connectors of the body device may be poor, or even non-existent, if the plug-in unit is not properly installed in its place. In one prior-art solution, pins of some poles in the electrical connectors between the plug-in unit and the body device are shorter than pins of other poles in these electrical connectors. These shorter pins are part of a test circuit for detecting whether the plug-in unit is installed in such a way that the shorter pins make galvanic contacts in the electrical connectors. If the shorter pins make the galvanic contacts, there is a good certainty that also the longer pins in the electrical connectors make adequate galvanic contacts. The above-described prior-art solution for indicating correct installation of a plug-in unit is suitable in cases where some poles of an electrical connector can be made such that galvanic contacts between these poles constitutes a sufficient guarantee of correct galvanic contacts between other poles of the electrical connector. In conjunction with some electrical connector structures this may be, however, challenging. Therefore, there is a need for other technical solutions for indicating correct installation of a plug-in unit."}
{"text": "1. Field of the Invention\nThis invention relates to a method for classifying and counting white blood cells, erythroblasts, and bacteria in body fluids by means of an automated hematology analyzer. More particularly, this invention relates to a method for simultaneously differentiating and counting white blood cell populations, erythroblasts, and bacteria in body fluids by means of multi-angle light scatter, fluorescence, and triple triggering circuitry in a three-dimensional space.\n2. Discussion of the Art\nExamination of various body fluids is critical for the diagnosis of bacterial meningitis, bacterial pneumonia or lung abscess, infection of the peritoneal cavity, and septic arthritis. The conventional method of analyzing body fluids in order to determine the presence of a bacterial infection, which involves dilution of biological samples, counting cells by means of a hemocytometer, preparing cell cultures, Gram staining, and microscopic examination, is tedious, time-consuming, and labor-intensive, and some clinical cases, such as bacterial meningitis, require immediate treatment because an untreated case can be lethal. Thus, the ability to analyze body fluids on a rapid hematology analyzer would be extremely useful.\nAnalysis of most body fluids drawn from hospitalized patients must be carried out in the hospital as soon as possible because such body fluids are not very stable and can be expected to deteriorate within approximately two hours. Cerebrospinal fluid can be expected to deteriorate within one hour. See, for example, Body Fluids: Laboratory examination of amniotic, cerebrospinal, seminal, serous, & synovial fluids: a text book atlas/C. Kjeldsberg and J. Knight, eds. 3rd ed. ASCP Press, 1993, incorporated herein by reference. Thus, analyzing body fluids rapidly on an automated hematology analyzer would be desirable in hospital laboratories.\nA number of manufacturers of hematology analyzers have systems that use the analysis of body fluids for cell counting. The Beckman-Coulter LH 750 hematology analyzer uses VCS technology (Volume by impedance, Complexity by radio-frequency, and laser light Scatter) for analysis of white blood cell differential. However, VCS technology cannot discriminate signals generated by bacterial cells from signals generated by other cell debris. The Bayer ADVIA® 2120 hematology analyzer uses myelo-peroxidase staining and light scatter to differentiate white blood cells. In the basophil channel, also known as the Lobularity/Nuclear density channel, a hypotonic surfactant solution is used to strip the cytoplasmic membrane from all leukocytes, except basophils. Neither the myeloperoxidase channel nor the basophil channel of the ADVIA® 2120 hematology analyzer is capable of distinguishing the signals generated by bacteria from signals generated by erythroblast nuclei or other cell debris.\nThe Sysmex XE-2100 hematology analyzer uses forward light scatter and side light scatter for counting white blood cells and nuclear staining, and side light scatter and fluorescence for differential analysis. However, the analyzer cannot distinguish the small noise signals generated by cell debris from those generated by bacteria. U.S. Pat. No. 5,325,168 describes a method and apparatus for analyzing cells in urine using both light scatter for determining size and fluorescence for determining differential DNA-staining intensity. This patent does not disclose how signals generated by small bacteria can be distinguished from noise signals generated by cell debris or from erythroblast nuclei. The cytograms of light scatter vs. fluorescence, i.e., FIGS. 14A, 14B, and 14C of U.S. Pat. No. 5,325,168, show no noticeable separation of noise signals from small bacterial signals.\nTo resolve the problems stated above, a rapid analysis of body fluids by means of an automated hematology analyzer, available in most clinical laboratories, is highly desirable to save the lives of infected patients by the rapid diagnosis of the medical condition of the patients and the subsequent treatment of the patients."}
{"text": "In current cellular or wireless technologies, the system controls are dominated by downlink (DL) measurements (UE measure and report). For example, the clustering or grouping of transmission points (TPs) for downlink (DL) and uplink (UL) transmissions is typically established based on user equipment (UE) measurements from DL reference signals, e.g., from each cell. This requires each cell to provide dedicated DL reference signals, which leads to complicated DL reference signal design and significant signaling overhead, and other issues for current or future wireless network technologies, such as in heterogeneous network (HetNet) where Pico cells may not provide reference signals. Further, using DL-based measurement, UL coordinated multi-point (CoMP) transmission may experience timing advance (the length of time a signal takes to reach TPs from a UE) issues within serving TPs. Another issue is that a dynamic (e.g., upon demand) TP optimization scheme requires frequent UE measurement and feedback, which is challenging for DL pilot design with more antennas. Moreover, the current UL sounding reference signal (SRS) is coupled with cell/TP identity and set up after TP clustering has been completed. Thus, mobile users have to change sounding signal configurations when switching cells. There is a need for a scheme to resolve or handle the issues above."}
{"text": "1. Field of the Invention\nThe invention is based on a gear pump, in particular for pumping fuel from a fuel tank to a high-pressure fuel pump.\n2. Description of the Prior Art\nOne such gear pump is known from German Patent Disclosure DE 101 34 622 A1. This gear pump has a housing, in which a pump chamber is formed in which a pair of gear wheels meshing with each other on their outer circumference is disposed. The gear wheels pump the medium to be pumped along pumping conduits, formed between their outer circumference and adjacent circumferential walls of the pump chamber, from a suction chamber into a pressure chamber. In operation of the gear pump, a tooth gap of one gear wheel reaches the pressure chamber, and a tooth gap of the other gear wheel reaches the pressure chamber after it, in chronological succession. The angle of rotation of the gear wheels corresponding to this spacing amounts here to 360°/2z in each case, where z is the number of teeth of the gear wheels.\nUpon the entry of a tooth gap into the pressure chamber, the volume of the pressure chamber is increased, causing the pressure in the pressure chamber to drop. In operation of the gear pump, pressure pulsations thus occur in the pressure chamber, which have a period corresponding to half the tooth pitch angle of the gear wheels. For a number of teeth z=10, the period is thus 18°, relative to the respective axes of rotation of the gear wheels. In other words, after 18° of an angle of rotation of each of the gear wheels, one tooth gap enters the pressure chamber, and the pressure chamber volume is increased. The pressure pulsations lead to high mechanical stress on the gear pump, which must be taken into account by means of an expensive construction of the gear pump that has the requisite strength. Moreover, the pressure pulsations reduce the pumping capacity of the gear pump. To reduce these pulsations, it has been proposed in DE 101 34 622 A1 that grooves be provided, by which the radial spacing between the outer circumference of the gear wheels and the circumferential walls is increased. By means of this provision, it has already been possible to achieve a significant reduction in the amplitudes of the pressure pulsations."}
{"text": "A variety of biological fluid collection devices and systems have been developed to assist with or otherwise aid in the production and collection of biological fluids. Included are systems and devices that aid in the stimulation and therefore production of seminal fluid such as masturbation sleeves and other such systems. Examples of such devices are disclosed in: U.S. Pat. No. 5,782,818, issued Jul. 21, 1998, titled “Device for Discreet Sperm Collection”; U.S. Pat. No. 5,806,523, issued Sep. 15, 1998, titled “Prophylactic and Prosthetic Device”; U.S. Pat. No. 5,807,360, issued Sep. 15, 1998, titled “Device for Discreet Sperm Collection”; and U.S. Pat. No. 5,466,235, issued Nov. 14, 1995, titled “Female Functional Mannequin”; all to Steven A. Shubin, Sr. and all of which are incorporated herein by reference in their entirety.\nThe collection of seminal fluid, whether in humans or animals, can be problematic and difficult. Whether, as is the case with human males collection can be an embarrassing and awkward experience, or as can be the case with male animals, collection can be problematic when attempting to obtain the sample. Further, and in the cases of large animals such as horses and the like, the collection can be dangerous due to the size, unpredictability, and un-cooperation of the animal.\nWith respect to humans, the procedure is still quite archaic and comprises the individual masturbating into a specimen receptacle, which can lead to discomfort. That is, the erect male penis is generally upwardly oriented, while the collection device must be held or situated lower or downwardly thereto. As such, the donor must force the erect penis downwardly in order for the ejaculate to be received in the specimen container. Forcing the erect penis downwardly not only can be difficult and uncomfortable, but in some cases can also be painful. Further, conflict between the position of the specimen container with respect to the erect penis can then lead to loss of ejaculate and/or inadvertent loss of more sperm rich ejaculate due to improper timing of the onset and/or ejaculation of the male. Whether due to discomfort that may occur during use or the difficulty in the simultaneous angulation and alignment of the specimen container and the penis, improvements are desired.\nYet further, the specimen containers that are used in today's facilities are designed for utilitarian purposes rather than comfort. For example, the specimen container itself functions strictly as a receptacle and offers nothing in the way of comfort, ease of use, or other assistance to and/or for the donor. Still further, the size, shape, rough edges and screw threads of today's specimen containers can actually cause the donor discomfort and even pain if contact with the genitalia is made.\nYet further, the collection of seminal fluid in animals is also problematic and difficult. For these specimens, collection can pose physical dangers to personnel engaged in the collection. As such, devices, systems, and methods that reduce the unpredictability and increase the cooperation of the animal are desired. Further, as the acquired semen specimens are not only used in the diagnoses and treatment of the animal, but may also be sold to breeders and the like, such specimens can be profitable and thus, the owners of these animals seek to have the specimen fully captured with the least amount of waste.\nMoreover, doctors say prostate health may benefit from the stimulation of the production of seminal fluid and thus systems that facilitate the process while mitigating the disadvantages previously described are desired. Recent studies have found that prostate health in human males may be related to frequency of ejaculation. In particular, infrequent ejaculation can lead to swelling of the prostate, known as congestive prostatitis, and may also increase the cancer risk in human males. Some medical sources suggest an ejaculation frequency of three to four times per week to ensure good prostate health. One study found a 14% lower lifetime prostate cancer rate for men who ejaculate between 13 and 20 times per month, and a upwards of 33% lower lifetime prostate cancer risk for men who ejaculate 21 times or more each month. Devices for stimulation and collection of seminal fluids may aid achieving higher ejaculations rates among men, particularly the unmarried and long-married.\nBeyond assisting with respect to prostate health effects of ejaculation, devices for stimulation and collection of seminal fluids through ejaculation may also assist in reversing desensitization issues. That is, repeated masturbatory stimulation of the penis using the hand or rough cloth can lead to desensitization of the penis, particularly in the absence of lubrication in some form. Desensitization can then result in erectile dysfunction during copulation. Use of properly lubricated devices designed specifically for the stimulation and collection of seminal fluid may help reverse the desensitization issues, and thus reduce the occurrence of erective dysfunction related to the desensitization issue."}
{"text": "In the past few years the electronics industry has increased the use and manufactured quantity, as well as demand for higher quality, of high purity alumina substrate chips for circuitry substrates used in component parts of computers and other devices such as switches and calculators. These chips need to be produced from a very pure alpha alumina, utilized as a substrate for the deposition of the various electronic pathways. Impurities in this alumina would cause leakage of the electric current through the substrate and cause dielectric losses which make such chips unreliable. One method for producing alpha alumina consisted of the production of a sol-gel of alumina followed by drying and subsequent calcination to the alpha form. As existing, this method produces impure material with a very hard structure resulting in large expenditures of energy to reduce the particulate to the proper size."}
{"text": "White LED lighting, known as the fourth generation of lighting source, has the advantages of solidified, small size, low heat, low power consumption, long life, fast responding speed and environmental, and is expected to be widely applied in the fields of general lighting and backlight sources in the future. In particular, white LED lighting, since it satisfies the requirements of energy saving and environmental protection of the green lighting projects, it is expected to be widely used in the future. However, the light principle of semiconductor LEDs determines that the emitted light can only be monochromatic light, while visible sunlight is a colorful spectrum composed of the colors of red, orange, yellow, green, cyan, blue and purple. It is obvious that a semiconductor LED light source only having a monochromatic light is not suitable for the field of general lighting. If the semiconductor LED light source is required to use for the field of general lighting, the defect of monochrome light emitting must be overcome through the application of other technologies. White LED technology is a typical representative of the transformation of LED monochrome light; the purpose thereof is to make the LED light source close to the natural sunlight color, thereby making it useful in the field of general lighting.\nPresently, in the manufacture of white LEDs at home and abroad, fluorescent material is a very important technology. White light is generated almost completely by using yellow phosphor combined with blue light. The property and preparation process of the fluorescent material directly affects the luminous efficiency, conversion efficiency, color coordination, color temperature and color-rendering property of the white LED. In actual application, the fluorescent material of the current white LED uses the amorphous phosphor as the main body. Generating a white light by the combination of single chip and phosphor is still the mainstream of the development of the general white light LED product. However, there are still several technical problems to be solved in the process of producing LEDs by using phosphor technology at the present stage of white LED: 1) the excitation efficiency and light conversion efficiency of phosphor are low; 2) the phosphor particles and dispersal uniformity issues are difficult to be effectively and completely solved; 3) phosphor lacks the red light-emitting component, it is difficult to produce a white LED with low color temperature and high color-rendering index; 4) the high light attenuation of phosphor leads to a shortened white LED life; and 5) phosphor has a poor physicochemical property that is unsuitable for the development requirement of the high-power LED. Under the influence of these problems of the properties of phosphor, the solving of the problems of white LED such as luminous efficiency, color-rendering property, life, high-power application and the like have encountered technical bottlenecks. This situation shows that the bottleneck of LED development increasingly highlighted that phosphor cannot meet the requirements of existing white LED and adapt to the future development trend of LED fluorescent lamp.\nIn the existing mixed LED white light technology, the white light is generated by a mixed light made by two or three complementary two-color LED light-emitting diodes or three primary color light-emitting diodes (LEDs). This technology generally uses only a simple mixing structure to obtain the mixed light, its color-rendering property is poor, and it is difficult to meet the practical requirements.\nMeanwhile, with the acceleration of the speed of white LED lights coming to the market and the expansion of the application field, it exposed a major application problem of “visual defect” that still existed in the current white LED lights. This defect bears the brunt of the harm and influence to the human eye, attracting the attention and research of the customer and all sectors of the society.\nIn view of the above problems, the present invention is aimed to design a solar spectrum-like LED structure. A solar spectrum can be simulated by gathering and mixing LED light of different colors. The solar spectrum-like LED structure of the present invention overcomes the technical defect that the traditional white LED technology needs to use phosphor and can simulate a solar spectrum, while, at the same time, has good color-rendering property and visual effect, and can be widely applied in the fields of general lighting, agriculture, animal husbandry and new biological energy sources."}
{"text": "During vascular surgery or endovascular treatment of vessels including thrombectomy, atherectomy, balloon angioplasty, and/or stent deployment, debris such as plaque and blood clots can move from the treatment site through a vein or artery and compromise the flow of blood at a location removed from the treatment site. In particular, various protection systems have been developed to prevent such debris from embolizing in the vessel. Distal protection devices include filters and occlusive devices (e.g., balloons) placed distally of the treatment site. Proximal protection devices include filters and occlusive devices placed proximally of the treatment site. In the case of filters, emboli collect within or on the filter. The filter with captured emboli is typically collapsed into a recovery catheter and the catheter withdrawn from the patient's body.\nThe size or number of emboli that must be retained by the filter in order to prevent clinically undesirable sequaelae is unknown. This uncertainty adds to the complexity of designing a filter with the appropriate characteristics. Small particles might pass through the filter pores and lodge downstream in tissues where they may cause tissue ischemia or tissue necrosis. In the heart, blood can be drawn and measurements can be made to track enzyme levels and determine myocardial damage. However, in the brain there is no easy and inexpensive method to evaluate the effect of a shower of emboli. Within the downstream tissue bed, there is a statistical component to the consequences of an embolus. For example, a 100 micron particle may lodge in a part of the brain where few adverse consequences are detected clinically, or it can lodge in a retinal artery, resulting in blindness in one eye. Therefore, it may be necessary to adjust the filter characteristics to suit the region of emboli filtration. A smaller pore size filter may be needed if protecting the brain than protecting the heart or kidney.\nEmbolic protection filters permit the passage of blood while retaining emboli that are larger than the pore size of the filter. Filter meshes are commonly made by incorporating holes in a polymer film, by interweaving filaments, or by producing interconnected porosity in a sheet of material (e.g., foam). It is difficult to make an embolic protection filter with the appropriate combination of pore size, pore area, embolic capacity, patency, mechanical strength, low collapsed or retracted profile, and recovery characteristics. Embolic filters made from polymer films commonly have a narrow range of pore sizes but suffer from a low percent open area because there is a limit to how closely the holes can be placed. Too little spacing between holes can result in a weak film that tears upon filter recovery. Foams tend to be bulky, thereby compromising the collapsed profile, and they have low strength.\nInterwoven meshes such as braids have the advantage of a pore area which is a high percentage of the total mesh area, excellent strength, and good flexibility, but tend to be made and used in ways that result in a wide range of pore sizes. A wide range of pore sizes is undesirable for a number of reasons. Patency is influenced by pore size. Theoretically, blood can be sheared as it flows through the pore, particularly at the edges of the pore opening. Shearing of blood can activate platelets and initiate a cascade of events that cause blood clotting. When filters are used in the bloodstream, it is common for thrombus to form in the vicinity of the smallest pores and no thrombus to form in the vicinity of the largest pores. Flow through the filter is thereby reduced because part of the filter becomes occluded. In addition, while some filters have a reasonable average pore size, a wide range of pore sizes in these filters may allow large particles to pass through the large pores during either the capture or recovery phase.\nA need in the art remains for an embolic protection filter having pores which are both small in size and which do not vary in size beyond an acceptable range."}
{"text": "Increasingly, wireless communications devices such as wireless telephones are becoming adapted for use with the Internet. It is estimated that, by the year 2003, half of all Internet traffic will come from wireless devices. Many present-day wireless telephones have at least some capability not only to capture and render audio information, but also to allow users to interact with data using a visual display and some form of data input device. Many wireless carriers are marketing access to the so-called “wireless web” as part of wireless telephone service.\nWhile wireless data access is clearly a boon to electronic device users, there is a notable deficiency in the current wireless data access technology: a user must generally choose to interact with the data either using the voice components of the wireless device (i.e., microphone and speaker), or using the visual components (i.e., screen and keypad), but cannot use both at the same time. Some aspects of communication work best with a visual interface, and others work best with a voice interface. For example, suppose that an application provides directions to a specified location—e.g., for use while driving. It is convenient to speak the name of the desired location as input to the application, but it is cumbersome to receive the directions themselves in the form of speech. A visual map combined with written directions such as “turn right on Elm” is a very convenient format in which to receive the directions, but it is may be less convenient to input the desired location using a keypad or stylus than it is merely to speak the location into a microphone. Thus, the ideal interface for certain applications is, in many cases, not visual or voice alone, but rather a combination of visual and voice. Present wireless application often allow one or the other but not both.\nOne problem that impedes the integrated and combined use of voice and visual interfaces to data is that each mode of communication generally requires its own browser. Typically, a particular piece of content (e.g., a web page) may be represented in both a visual markup language (such as Wireless Markup Language or “WML”), and in a voice markup language (such as Voice eXtensible Markup Language or “VXML”). A visual browser permits the user to navigate through the WML content using the screen and keypad. Similarly, a voice browser, which is generally a software component separate from the visual browser, permits the user to navigate through the VXML content using the microphone and speaker. Not only are the visual and voice browsers separate software components; they often execute on separate, and distantly located, devices. A visual browser typically executes on a wireless handset (such as a wireless-web-enabled telephone). However, the handset is generally “dumb” with respect to voice—i.e., it can capture and render audio signals, but does not have the means to browse and navigate content based on the content of the received audio, or to generate audio signals based on VXML data. Thus, a voice browser typically executes on a voice server and communicates with the user through the microphone and speaker of the wireless device by transmitting and receiving digital signals to the device through an ordinary voice circuit within a wireless network.\nBecause the voice and visual browsers are separate, it is difficult to switch seamlessly back and forth between a visual and voice modes of interacting with wireless data, because the voice browser may be unaware of what the visual browser is doing and vice versa. That is, the voice and visual browsers are not normally “synchronized,” in the sense that neither knows the other's state with respect to the underlying content that the voice and visual browsers are manipulating. For example, suppose that a wireless handset user uses a visual browser to navigate through a series of web pages, eventually ending up at a particular URL. If the user then decides to switch to the voice interface, the voice browser does not know where the user has navigated to because it is unaware of what the visual browser has been doing. Upon switching to a voice interface, the voice browser can simply re-start the user at a “home” URL, but this is inconvenient for the user because the user loses the benefit of all of the navigation that has already been performed. Similarly, within the page located at a particular URL, the user may have navigated through several cards, and may have positioned the cursor at a particular field on a particular card using the visual browser, but the voice browser will be unaware of all of this activity. The problem, in this example, is that the voice and visual browsers are not “synchronized.”\nIn view of the foregoing, there is a need for a browser synchronization technique that overcomes the drawbacks of the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to an information processing apparatus that can manage a peripheral device, a peripheral device control method, and a program.\n2. Description of the Related Art\nA peripheral device control system enables a user of an information processing apparatus to access a peripheral device via an appropriate interface, such as Universal Serial Bus (which may be referred to as USB in the following description), Ethernet®, or wireless LAN. This kind of control systems can be effectively used for various users in their houses and offices. An example of the peripheral apparatus is, for example, a printer, a copying machine, a facsimile machine, a scanner, a digital camera, or a multifunction peripheral.\nFurther, Windows®7 (hereinafter, referred to as W7), which is provided by Microsoft Corporation, includes newly introduced functions to manage peripheral apparatuses connected to a personal computer (hereinafter, referred to as PC) or comparable information processing apparatus.\nIn W7, Device Stage® can be thought of as a home page for peripheral hardware such that when a new peripheral device is connected to a PC (device running W7), a menu of tasks for the newly connected peripheral is shown in one place. Device Stage includes a “Devices and Printers folder” serving as a window that displays a peripheral device connected to the PC. Further, Device Stage® can display a link to a unique application or service provided by each peripheral device. Hereinafter, the “Devices and Printers folder” is referred to as “DP folder.”\nThus, it becomes feasible to provide a task (button) on a Device Stage screen as a link to a function for displaying an electronic manual (EM) dedicated to each peripheral device. In this case, the Device Stage screen enables users to launch (execute) a manual viewer application to display an EM for an intended peripheral device.\nAn example of the peripheral device is a multi function peripheral (MFP) device, such as a printer with capability to scan, fax and email documents. Further, an example of the application is the manual viewer application. The manual viewer application can be used to display an EM for each peripheral device. In general, the above-described application can support a plurality of peripheral devices and a plurality of languages.\nThe application can display a user interface that includes a peripheral device selection field and a language selection field to enable users to select a peripheral device or select a language. Thus, each user can use a preferred language to display and execute functions of a target peripheral device.\nThe electronic manuals, which can be displayed by the manual viewer application, are prepared for each peripheral device using a plurality of languages with reference to a plurality of destinations thereof. A setup program determines a language of the EM for each peripheral device, when the EM is installed, with reference to a view language of the Operating System (which may be referred to as OS in the following description). Further, users may be allowed to select a preferred version (language) of the EM for each peripheral device.\nIf an application does not have any function for selecting a peripheral device or a language, the application constantly uses the peripheral device and the language that have been selected or designated in an installation process.\nFurther, if an application has a function for selecting a peripheral device or a language, the application preliminarily sets a peripheral device and a language as initial values selected or designated in an installation process. In this case, after the installation is completed, users are allowed to change the setting in a selectable range so as to designate an intended peripheral device or a preferred language.\nFurther, the information processing apparatus capable of performing the above-described processing is, for example, discussed in Japanese Patent Application Laid-Open No. 2005-85132.\nAs described above, a static text file having an Extensible Markup Language (XML) format cannot include any executable script. Therefore, performing file generation, file editing (especially, writing), generation of registry information storing setting information, and editing (especially, writing) of the registry information is difficult.\nAccordingly, when a user executes a function (presses a button) on the Device Stage (peripheral device management screen), it is difficult to designate an appropriate language set by the user and execute (display) the function."}
{"text": "Managing large document databases may be a daunting task. While a user may wish to browse the whole database, a user also may wish that documents are sorted with respect to visual similarity and/or layout features. Simply presenting documents in a random order, however, is less than ideal."}
{"text": "Exhaust gas recirculation is a technique commonly used for controlling the generation of undesirable pollutant gases and particulate matter in the operation of internal combustion engines. This technique has proven particularly useful in internal combustion engines used in motor vehicles such as passenger cars, light duty trucks, and other on-road motor equipment. The exhaust gas recirculation technique primarily involves the recirculation of exhaust gas by-products into the intake air supply of the internal combustion engine. This exhaust gas thus reintroduced to the engine cylinder reduces the concentration of oxygen therein, which in turn lowers the maximum combustion temperature within the cylinder and slows the chemical reaction of the combustion process, decreasing the formation of nitrous oxide. Furthermore, the exhaust gases typically contain a portion of unburned hydrocarbon which is burned on its reintroduction into the engine cylinder, which further reduces the emission of exhaust gas by-products which would be emitted as undesirable pollutants from the internal combustion engine.\nAnother technique useful in the control and reduction of undesirable emissions from internal combustion engines is the use of pressure-charged intake air. This permits the use of relatively smaller cubic displacement and lighter weight internal combustion engines in mobile equipment, reducing in turn the specific fuel consumption of the vehicle and overall mass of the vehicle necessary to perform a given function. In addition to the benefits of reduced size and mass, the typical pressure-charging device may be controlled to provide improved emissions characteristics. Pressure-charging machines suitable for such applications include the exhaust gas driven turbocharger which is comprised typically of an exhaust gas driven turbine linked to a turbine disposed in the intake air stream to provide compression of the intake air. The typical turbocharger is controlled by providing a gate which controls exhaust gas flow and gates exhaust gas to bypass the exhaust gas turbine and control the charging rate of the turbocharger so that the maximum pressure limits of the associated internal combustion engine are not exceeded.\nStill another technique for controlling emissions used by many engine manufactures is the use of aftercooling the intake air thereby reducing the intake manifold temperature. Some of the related art techniques have also considered intercooling the EGR gases by routing the recirculated exhaust gases through an aftercooler. It is well known that lower intake manifold temperatures tends to reduce the formation of nitrous oxides found in the exhaust.\nWhen utilizing EGR in a turbocharged diesel engine, the exhaust gas to be recirculated is preferably removed upstream of the exhaust gas driven turbine associated with the turbocharger. In many EGR applications, the exhaust gas is diverted directly from the exhaust manifold. Likewise, the recirculated exhaust gas is preferably re-introduced to the intake air stream downstream of the compressor and air-to-air aftercooler. For example, in many EGR applications the recirculated exhaust gas is reintroduced to the intake manifold.\nReintroducing the exhaust gas downstream of the compressor and air-to-air aftercooler is preferred due to the reliability and maintainability concerns that arise should the exhaust gas is passed through the compressor and aftercooler. However at some engine operating conditions, there is a pressure differential between the intake manifold and the exhaust manifold which essentially prevents many conventional EGR systems from being utilized. For example, at high speed, high load conditions in a turbocharged engine, the exhaust gas does not readily flow from the exhaust manifold to the intake manifold. What is needed, therefore, is a simple and inexpensive technique for recirculating exhaust gas from the exhaust manifold to the intake manifold at all engine operating conditions.\nAnother problem associated with many conventional EGR systems is that the turbocharger efficiency is sacrificed when exhaust gas is diverted from the exhaust manifold. Removing the exhaust gas to be recirculated from the exhaust manifold or elsewhere upstream of the exhaust gas driven turbine depletes the mass flow and heat energy passing through the turbine which, in turn, lowers the boost levels created by the compressor. Most diesel engine turbochargers are fixed geometry turbochargers, in that they are specifically designed to operate efficiently when matched to the engine exhaust flow output. The reduction in mass flow and pressure due to the EGR creates a mismatch between the exhaust flow to the turbocharger and the turbine specifications during EGR operation. The mismatch results in a turbocharger output that is reduced in percentage more than the percentage reduction in exhaust flow to the turbocharger thereby creating significant losses in airflow and boost pressure. The reduction in airflow and boost pressure decreases the air to fuel ratio down to a point where particulates as well as the brake specific fuel consumption (BSFC) increase. Disadvantageously, the reduction in airflow and boost pressure also results in a noticeable difference in engine performance to the operator depending on whether EGR is on or off."}
{"text": "The present invention relates generally to theft prevention devices and, more particularly, to a security device that is specifically adapted to secure a laptop computer to a stationary or portable support or stand such as a technology workstand for securing the laptop computer in the open position so as to effectively prevent theft of the laptop computer.\nComputers and computer-related devices have undergone tremendous advancement in recent years. Evidence of this rapid technological growth can be seen in the increasing power of computers in concert with increasing miniaturization of computers and peripherals. For example, computer peripherals such as display monitors are commercially available with increasing amounts of features but in greatly reduced size as evidenced by thin, flat panel monitors that are replacing traditional display monitors.\nAlthough such improvements in computer technology are invaluable in terms of their operational capabilities, certain deficiencies present issues which detract from the overall utility of such devices. In particular, the ultimate reduction in size of computers and computer peripherals has made these items easier to carry and/or conceal. Obviously, such reduction in size has drastically increased susceptibility to theft.\nDue to the relatively high expense of computers and peripherals, financial losses due to theft is of special concern for entities such as corporations, government agencies and educational institutions since such entities typically purchase large quantities of these items. As such, it is of primary importance for such entities to address how theft can be reduced prior to making computers and associated devices available to their employees, students and others who may use these computers in public and semi-public areas where theft is of great concern.\nIn order to reduce losses from theft, various types of anti-theft computer locks are currently available from manufacturers of security products. Although such currently known anti-theft locks may achieve their primary objective of computer protection, they possess certain deficiencies which fail to optimally protect the new generation of computers and computer-related devices.\nPerhaps the greatest deficiency associated with conventional anti-theft computer locks is their inability to secure a variety of laptops or other similar devices. For example, certain computer locks must engage and secure a specifically shaped micro-security slot that may be provided by certain computer manufacturers. This limitation has characterized many anti-theft computer locks of the prior art. In this regard, the technological innovation in the field of computer locks has been minimal compared to overall advances in computer technology.\nA further deficiency associated with existing computer locks relates to the cables used in certain anti-theft computer locks. Generally, these cables are fixedly attached to and extend from the computer lock at portions which are opposite from their slot engaging portions. Due to the fixed attachment of the cables to the computer locks, the cable portions near such attachment point may become locally strained if secured to a stationary object that is placed in angular orientation to the computer lock. The lock also requires the cable to loop around a stationary object to prevent theft of the laptop computer. Thus, transporting the laptop computer from one room to another room requires unlocking the computer lock to remove the cable from the stationary object in one room, and then re-locking the computer lock after looping the cable around a stationary object in another room.\nIn view of the above-described shortcomings of conventional anti-theft computer locks, there exists a need in the art for a security device for securing a laptop computer to a fixed surface in an effective manner. More specifically, there exists a need in the art for a security device that can effectively secure a laptop computer to a support panel and which further provides the capability of protecting the device against theft when the laptop computer is oriented in an open position wherein the display monitor is pivoted relative to the laptop keyboard."}
{"text": "Starch is an industrial and food raw material of major importance. Starch for food use is commonly extracted from starchy plant constituents (cereal grains, seeds, tubers) using procedures that involve reducing the plant material to fine particles, washing the starch from the insoluble plant material as a suspension in water, then removing the starch from the suspension, usually by a process involving filtering or centrifugation in some form. To obtain starch of high purity, it is important that the non-starch plant materials have, or can be made to have, differences in solubility or density, compared to the starch.\nThe process of separating starch from insoluble and soluble components of plants is an important step for obtaining purified starch, but it is also important for the isolation of other plant components that may be valuable, where removal of the starch is a required step.\nIndustrial extraction of starch using procedures based on the above principle is applied in various forms. For example, starch is isolated from potatoes and cassava by rasping the raw material to open the plant cells, reducing the solubility of proteins by reaction with sulphur dioxide or sodium bisulphate, then washing and extracting the starch. Starch is isolated from maize (zea mays) by steeping then wet milling the corn to remove the germ and husk, fine grinding, then applying centrifugal separation and filtering techniques to separate the starch from the proteins. Starch is isolated from wheat by mixing milled wheat flour with water to a dough then washing the starch from the insoluble wheat gluten.\nStarch may also be isolated from plant materials by non-aqueous extractions. For example, in a review of procedures for obtaining starch from barley, McDonald and Stark (1) identified several procedures, but considered none food safe or industrially applicable.\nThe aqueous separation processes described above are applicable to a limited number of plants, and thus most industrial starch is made from only four sources: potatoes, cassava, maize and wheat. Rice starch can be isolated from rice using a process involving steeping in alkaline solution, but with low efficiency. The main barrier to extraction of starch from other plant sources is the difficulty in reducing the solubility of large non-starch plant polymers sufficiently to allow separation from the starch."}
{"text": "The present invention relates to a new and distinct cultivar of Hydrangea plant, botanically known as Hydrangea macrophylla and hereinafter referred to by the name ‘Hokomarevo’.\nThe new Hydrangea plant is a naturally-occurring whole plant mutation of Hydrangea macrophylla ‘Xian’, not patented. The new Hydrangea plant was discovered and selected by the Inventors as a single flowering plant from within a population of plants of ‘Man’ in a controlled greenhouse environment in Boskoop, The Netherlands in July, 2007.\nAsexual reproduction of the new Hydrangea plant by vegetative cuttings in a controlled greenhouse environment in Boskoop, The Netherlands since August, 2007, has shown that the unique features of this new Hydrangea plant are stable and reproduced true to type in successive generations."}
{"text": "In a high-speed transmission and reception module for optical communication, for the connection between an optical device and a circuit substrate, a flexible substrate manufactured using a transmission line has heretofore been used for correcting a displacement between the optical device and the circuit substrate.\nHowever, the above-described high-speed transmission and reception module for optical communication is hard to realize sufficient high-frequency characteristics at a communication speed more than 10 Gbps. To cope with the above-described problem, in the high-speed transmission and reception module for optical communication having a communication speed more than 10 Gbps, a relay substrate made of ceramic has been connected between the flexible substrate and the circuit substrate as a substrate for relay.\nProposed is a connector element for high frequency transmission in which parts having surfaces facing to each other are provided on conductors for constituting a ground line and signal line of the connector element and covered with dielectric materials (see, e.g., Japanese Laid-open Patent Publication No. 06-215819).\nHowever, the above-described relay substrate made of ceramic is expensive. Further, precision is necessary for the connection between the relay substrate and any one of the flexible substrate and the circuit substrate, and also the time for assembly is necessary. Therefore, when the high-speed transmission and reception module for optical communication is manufactured, the assembly causes an increase in cost.\nIn the above-described connector element for high-frequency transmission, the connection shape to the substrate of a signal line and ground line constituting the transmission line is not sufficiently taken into consideration. Therefore, even if high-frequency transmission characteristics of the connector element are preferable, loss of high-speed signal may be caused by the connection to the substrate."}
{"text": "Noroviruses are non-cultivatable human Caliciviruses that have emerged as the single most important cause of epidemic outbreaks of nonbacterial gastroenteritis (Glass et al., 2000; Hardy et al., 1999). The clinical significance of Noroviruses was under-appreciated prior to the development of sensitive molecular diagnostic assays. The cloning of the prototype genogroup I Norwalk virus (NV) genome and the production of virus-like particles (VLPs) from a recombinant Baculovirus expression system led to the development of assays that revealed widespread Norovirus infections (Jiang et al. 1990; 1992).\nNoroviruses are single-stranded, positive sense RNA viruses that contain a non-segmented RNA genome. The viral genome encodes three open reading frames, of which the latter two specify the production of the major capsid protein and a minor structural protein, respectively (Glass et al. 2000). When expressed at high levels in eukaryotic expression systems, the capsid protein of NV, and certain other Noroviruses, self-assembles into VLPs that structurally mimic native Norovirus virions. When viewed by transmission electron microscopy, the VLPs are morphologically indistinguishable from infectious virions isolated from human stool samples.\nImmune responses to Noroviruses are complex, and the correlates of protection are just now being elucidated. Human volunteer studies performed with native virus demonstrated that mucosally-derived memory immune responses provided short-term protection from infection and suggested that vaccine-mediated protection is feasible (Lindesmith et al. 2003; Parrino et al. 1997; Wyatt et al., 1974).\nAlthough Norovirus cannot be cultivated in vitro, due to the availability of VLPs and their ability to be produced in large quantities, considerable progress has been made in defining the antigenic and structural topography of the Norovirus capsid. VLPs preserve the authentic confirmation of the viral capsid protein while lacking the infectious genetic material. Consequently, VLPs mimic the functional interactions of the virus with cellular receptors, thereby eliciting an appropriate host immune response while lacking the ability to reproduce or cause infection. In conjunction with the NIH, Baylor College of Medicine studied the humoral, mucosal and cellular immune responses to NV VLPs in human volunteers in an academic, investigator-sponsored Phase I clinical trial. Orally administered VLPs were safe and immunogenic in healthy adults (Ball et al. 1999; Tacket et al. 2003). At other academic centers, preclinical experiments in animal models have demonstrated enhancement of immune responses to VLPs when administered intranasally with bacterial exotoxin adjuvants (Guerrero et al. 2001; Nicollier-Jamot et al. 2004; Periwal et al. 2003). Collectively, these data suggest that a vaccine consisting of properly formulated VLPs represents a viable strategy to immunize against Norovirus infection."}
{"text": "High energy feed supplements are currently extensively employed in the raising of cattle, sheep, and other domestic animals. These feed supplements include such materials as molasses, fat, and urea, and these materials are usually sold separately and mixed by the animal feeder with such feeds as grains or mixed feeds.\nFor more convenient handling and distribution, various high energy feed supplement materials have heretofore been combined in two distinctive composite forms, namely, in the form of a liquid feed supplement, and in the form of a block type of feed supplement. However, both of these prior art types of composite feed supplement preparations have involved serious problems in production, handling and usage, and have been undesirably limited in the amounts of certain feed supplement substances which could be contained therein.\nThus, prior art liquid feed supplements have been available which contain a large percentage of molasses for high energy content, and which also contain a substantial amount of urea as a source of nitrogen for protein equivalent. However, the amount of fat content which could be provided in such prior art liquid feed supplements is severely limited, and is much less than is desirable for an optimum feed supplement composition. The biggest difficulties in connection with prior art liquid feed supplements have been in the handling and ultimate use thereof. Thus, special tanks are required for storage of liquid feed supplements, and tank trucks are required for primary distribution thereof. Additionally, putting such liquid feed supplements out for consumption by cattle or other domestic animals presents a very serious problem, as open troughs of the sweet, sticky liquids attract flies and other undesirable insects. Accordingly, special closed liquid feed dispensing containers have been required which meter the liquid feed supplement material out through mechanical devices. Even when thus dispensed, cattle tend to get the sweet, sticky feed supplement liquid all over their faces, which draws flies and is accordingly unsanitary.\nBlock types of animal feed supplements have heretofore been produced by combining chopped alfalfa, straw, or grains, and proteins, with from about 10% to about 30% molasses in a mixer, and then pressure forming the combination under extremely high pressure, as for example about 100 tons per square inch. This prior art method of forming block feed supplements involves the use of very costly machinery, and labor is also quite costly as it is principally a hand operation. In such prior art feed supplement blocks the maximum percentage of molasses which can be satisfactorily combined in the mixture is only about 3:%, as combinations with larger percentages of molasses are too much of a gooey mess for satisfactory pressure forming. Similarly, the amounts of fat which could be contained in such pressure formed feed supplements were inadequate.\nIn general, it is desirable in a feed supplement for animals to include a relatively large percentage of molasses, because of the high content of inverted sugars therein (over 50%), because of the good mineral content, and in particular because of the large amount of energy available cheaply. A substantial amount of fat has also been found to be desirable in such feed supplements. While prior art liquid feed supplements were able to obtain an adequate amount of molasses, they were generally deficient in the desirable amount of fat, and were difficult to handle and use. On the other hand, while prior art block type feed supplements were much more convenient to use, they were deficient in both the amount of molasses they could contain and the amount of fats they could contain, and they have the further disadvantage of undesirably high production cost and investment in machinery.\nPrior to the method of the present invention wherein the molasses is first heated at substantially ambient pressure to remove a major portion of the water content thereof and then without further heating subjected to a vacuum to further reduce the water content without foaming, there has been no method in the art capable of producing dense, hard, vitreous molasses feed blocks, and particularly such blocks wherein substantial quantities of other important animal feed ingredients are integrally combined in such blocks with the molasses solids serving as a supporting matrix.\nMolasses has heretofore been dried by vacuum methods for animal feed purposes. However, prior art methods of this type were directed to the drying of the molasses for later mixing with other feed ingredients, and did not produce a dense, hard, vitreous molasses feed block; nor were they capable of producing such a block. Examples of such prior art vacuum drying methods were the methods described in U.S. Pat. No. 1,983,434, issued Dec. 4, 1934 to Black et al, and U.S. Pat. No. 2,089,062 issued Aug. 3, 1937 to Hougland.\nIn the Black et al patent, prior to the removal of any water from the molasses, it was passed into a high speed vacuum evaporator at low temperature (150.degree. to 160.degree. F.), in which most of the water was removed, and then it was heated to between 200.degree. F. and 238.degree.F. and spray dried. The purpose of the Black et al method was to avoid having to mix the molasses and other feed ingredients while the molasses was in its liquid form, because the stickiness of the molasses made mixing so difficult.\nIn the Hougland patent, a thin film of molasses was applied to a heated surface in an evacuated drum, so that the vacuum was applied to the molasses prior to and during the heating thereof, and hence before the removal of any of the water from the molasses. The purpose of the Hougland method was to dry the molasses so that when it was later mixed with other feeds it would not be subject to fermentation.\nThe application of vacuum for drying before the removal by heating of any water from the molasses in the Black et al and Hougland methods precluded the formation of dense, hard, vitreous feed blocks by such methods; and also precluded the inclusion of other feed materials with the molasses when it was being treated, and hence precluded the integral inclusion of other feed materials within a matrix of hard molasses.\nSuch prior art methods also were rather severely limited in production rate because of the tendency of the molasses to swell and foam when thus subjected to the vacuum with initial water content, making animal feed production by such methods generally costly and inefficient."}
{"text": "The present invention relates to a magnet device for animals, in particular cattle, and particularly to a magnet device adapted to be arranged in the stomach of cattle and to collect scrap iron pieces taken down in the stomach so that the scrap iron pieces are prevented from entering into the intestines.\nCattle grazing in meadows sometimes swallow grass nails, wires or the like falling on the ground. These scrap iron pieces remaining in the stomach are slightly dangerous for cattle, but when they enter into the intestines from the stomach, it causes a great danger which can be fatal to the cattle.\nIn order to prevent such a danger, there has been proposed a device which consists of a magnet formed in a bar or cylinder at a size to not enter into the intestines from the stomach and which is adapted to remain in the stomach of cattle thereby collecting scrap iron or the like invading into the stomach.\nHowever, since such a device of a bar or cylinder magnet has magnetic pole portions limited to both end portions thereof, portion for performing a magnetizing function which collects the scrap irons is also limited to only both the end portions and its collecting effect has been insufficient.\nIn short, since the magnet device of this kind must remain in the stomach of cattle after the cattle has once swallowed it, without removing it out of the body of the cattle during its life, it was necessary to have the wide portion to magnetize scrap irons and to have a strong magnetizing force to magnetize the scrap irons.\nThe present invention is developed in the light of the above circumstances and in order to eliminate the above defects of the prior device, and contemplates to provide a magnet device which has the magnetizing portion not only at both ends but at peripheral wall portion and which has a magnetizing force intensified by concentrating flux of magnetic pole.\nAccording to the present invention, a magnetic device is constituted by a plurality of magnets in a short cylinder shape having magnetic poles of polarity differing from each other at both end surfaces and piled in a longitudinal direction through the intermediary of at least a magnetic plate between each magnet, magnetic poles of the magnet opposing through the intermediary of the magnetic plate being the same polarity whereby the flux concentrates at the magnetic plate (yoke) portion and the magnetizing force increases.\nFurther, according to the invention, a magnet device is constituted by a plurality of magnets in a short cylinder shape having those outer diameter which are similar to or smaller than the outer diameter of magnetic plate, in order to decrease leakage flux and to effectively utilize the flux.\nOther objects and features of the invention will appear in the following description taken with reference to the accompanying drawing."}
{"text": "Aerobic exercise is a popular form of exercise that improves one's cardiovascular health by reducing blood pressure and providing other benefits to the human body. Aerobic exercise generally involves low intensity physical exertion over a long duration of time. Generally, the human body can adequately supply enough oxygen to meet the body's demands at the intensity levels involved with aerobic exercise. Popular forms of aerobic exercise include running, jogging, swimming, and cycling among others activities. In contrast, anaerobic exercise often involves high intensity exercises over a short duration of time. Popular forms of anaerobic exercise include strength training and short distance running.\nMany choose to perform aerobic exercises indoors, such as in a gym or their home. Often, a user will use an aerobic exercise machine to have an aerobic workout indoors. One such type of aerobic exercise machine is stepping machine, which often includes foot supports that move along generally vertical arcuate paths when moved by the feet of a user. Other popular exercise machines that allow a user to perform aerobic exercises indoors include treadmills, rowing machines, elliptical trainers, and stationary bikes to name a few.\nOne type of stepping machine is disclosed in U.S. Patent Publication No. 2014/0274575 issued to Rasmey Yim, et al., (hereinafter “the '575 Publication”). In this reference, embodiments of stationary exercise machines are described as having reciprocating foot and/or hand members, such as foot pedals that move in a closed loop path. The '575 Publication, abstract. Some embodiments can include reciprocating foot pedals that cause a user's feet to move along a closed loop path that is substantially inclined, such that the foot motion simulates a climbing motion more than a flat walking or running motion. Id. Some embodiments are described as including reciprocating handles that are configured to move in coordination with the foot via a linkage to a crank wheel also coupled to the foot pedals. Id. Variable resistance can be provided via a rotating air-resistance based mechanism, via a magnetism based mechanism, and/or via other mechanisms, one or more of which can be rapidly adjustable while the user is using the machine. Id. According to this reference, traditional stationary exercise machines include stair climber-type machines and elliptical running-type machines. The '575 Publication, para. [0003]. Each of these types of machines typically offers a different type of workout, with stair climber-type machines providing for a lower frequency vertical climbing simulation, and with elliptical machines providing for a higher frequency horizontal running simulation. Id. Other types of exercise machines are disclosed in U.S. Pat. No. 5,242,343 to Miller; U.S. Pat. No. 5,499,956 to Miller; U.S. Pat. No. 5,540,637 to Rodgers; U.S. Pat. No. 5,573,480 to Rodgers; U.S. Pat. No. 5,683,333 to Rodgers; U.S. Pat. No. 5,938,567 to Rodgers; and U.S. Pat. No. 6,080,086 to Maresh. These references are incorporated herein by reference for all that they disclose."}
{"text": "Although fibers of the present invention are suitable for reinforcing various matrix materials, such as adhesives, asphalts, composites, plastics, rubbers, etc., and structures made from these, the fibers that will be described herein are especially suited for reinforcing hydratable cementitious compositions, such as ready-mix concrete, precast concrete, masonry concrete (mortar), shotcrete, bituminous concrete, gypsum compositions, gypsum- and/or Portland cement-based fireproofing compositions, and others.\nA major purpose of the fibers of the present invention is to reinforce concrete, e.g., ready-mix, shotcrete, etc., and structures made from these. Such matrix materials pose numerous challenges for those who design reinforcing fibers.\nConcrete is made using a hydratable cement binder, a fine aggregate (e.g., sand), and a coarse aggregate (e.g., small stones, gravel). A mortar is made using cement binder and fine aggregate. Concretes and mortars are hence brittle materials. If a mortar or concrete structure is subjected to stresses that exceed its maximum tensile strength, then cracks can be initiated and propagated therein. The ability of the cementitious structure to resist crack initiation and crack propagation can be understood with reference to the “strength” and “fracture toughness” of the material.\n“Strength” relates to the ability of a cement or concrete structure to resist crack initiation. In other words, strength is proportional to the maximum load sustainable by the structure without cracking and is a measure of the minimum load or stress (e.g., the “critical stress intensity factor”) required to initiate cracking in that structure.\nOn the other hand, “fracture toughness” relates to the specific “fracture energy” of a cement or concrete structure. This concept refers to the ability of the structure to resist propagation—or widening—of an existing crack in the structure. This toughness property is proportional to the energy required to propagate or widen the crack (or cracks). This property can be determined by simultaneously measuring the load required to deform or “deflect” a fiber-reinforced concrete (FRC) beam specimen at an opened crack and the amount or extent of deflection. The fracture toughness is therefore determined by dividing the area under a load deflection curve (generated from plotting the load against deflection of the FRC specimen) by its cross-sectional area.\nIn the cement and concrete arts, fibers have been designed to increase the strength and fracture toughness in reinforcing materials. Numerous fiber materials have been used for these purposes, such as steel, synthetic polymers (e.g., polyolefins), carbon, nylon, aramid, and glass. The use of steel fibers for reinforcing concrete structures remains popular due to the inherent strength of the metal. However, one of the concerns in steel fiber product design is to increase fiber “pull out” resistance because this increases the ability of the fiber to defeat crack propagation. In this connection, U.S. Pat. No. 3,953,953 of Marsden disclosed fibers having “J”-shaped ends for resisting pull-out from concrete. However, stiff fibers having physical deformities may cause entanglement problems that render the fibers difficult to handle and to disperse uniformly within a wet concrete mix. More recent designs, involving the use of “crimped” or “wave-like” polymer fibers, may have similar complications, depending on the stiffness of the fiber material employed.\nPolyolefin materials, such as polypropylene and polyethylene, have been used for reinforcing concrete and offer an economic advantage due to relative lower cost of the material. However, these polyolefinic materials, being hydrophobic in nature, resist the aqueous environment of wet concrete. Moreover, the higher amount of polyolefin fibers needed in concrete to approximate the strength and fracture toughness of steel fiber-reinforced concrete often leads to fiber clumping or “balling” and increased mixing time at the job site. This tendency to form fiber balls means that the desired fiber dosage is not achieved. The correct concentration of fibers is often not attained because the fiber balls are removed (when seen at the concrete surface) by workers intent on achieving a finished concrete surface. It is sometimes the case that locations within the cementitious structure are devoid of the reinforcing fibers entirely. The desired homogeneous fiber dispersion, consequently, is not obtained.\nMethods for facilitating dispersion of fibers in concrete are known. For example, U.S. Pat. No. 4,961,790 of Smith et al. disclosed the use of a water-soluble bag for introducing a large number of fibers into a wet mix. U.S. Pat. No. 5,224,774 of Valle et al. disclosed the use of non-water-soluble packaging that mechanically disintegrated upon mixing to avoid clumping and to achieve uniform dispersal of fibers within the concrete mix.\nThe dispersal of reinforcing fibers could also be achieved through packaging of smaller discrete amounts of fibers. For example, U.S. Pat. No. 5,807,458 of Sanders disclosed fibers that were bundled using a circumferential perimeter wrap. According to this patent, the continuity of the peripheral wrapping could be disrupted by agitation within the wet concrete mix, and the fibers could be released and dispersed in the mix.\nOn the other hand, World Patent Application No. WO 00/49211 of Leon (published Aug. 24, 2000) disclosed fibers “packeted” together but separable when mixed in concrete. A plurality of fibers were separably-bound together, such as by tape adhered to cut ends of the fibers, thereby forming a “packet.” Within a wet cementitious mix, the packets could be made to break and/or dissolve apart to permit separation and dispersal of individual fibers within the mix.\nThe dispersal of reinforcing fibers may also be achieved by altering fibers during mixing. For example, U.S. Pat. No. 5,993,537 of Trottier et al. disclosed fibers that progressively fibrillated upon agitation of the wet concrete mix. The fibers presented a “low initial surface area” to facilitate introducing fibers into the wet mix, and, upon agitation and under the grinding effect of aggregates in the mix, underwent “fibrillation,” which is the separation of the initial low-surface-area fibrous material into smaller, individual fibrils. It was believed that homogeneous fiber distribution, at higher addition rates, could thereby be attained.\nA novel approach was taught in U.S. Pat. No. 6,197,423 of Rieder et al., which disclosed mechanically-flattened fibers. For improved keying within concrete, fibers were flattened between opposed rollers to attain variable width and/or thickness dimensions and stress-fractures perceivable through microscope as discontinuities and irregular and random displacements of polymer on the surface of the individual fibers. This microscopic stress fracturing was believed to improve bonding between cement and fibers, and, because the stress-fractures were noncontinuous in nature, the fibers were softened to the point at which fiber-to-fiber entanglement (and hence fiber balling) was minimized or avoided. The mechanical-flattening method of Rieder et al. was different from the method disclosed in U.S. Pat. No. 5,298,071 of Vondran, wherein fibers were interground with cement clinker and retained cement particles embedded into the surface.\nIn this vein, the nature of the fiber surface has also been a frequent topic of research in fiber dispersion and bonding in concrete. For example, U.S. Pat. No. 5,753,368 of Hansen disclosed a list of wetting agents such as emulsifiers, detergents, and surfactants to render fiber surfaces more hydrophilic and thus more susceptible to mixing in wet concrete. On the other hand, U.S. Pat. No. 5,753,368 of Berke et al. taught that the bonding between concrete and fibers could be enhanced by employing particular glycol ether coatings instead of conventional wetting agents that tended to introduce unwanted air at the fiber/concrete interface.\nOf course, as mentioned in U.S. Pat. Nos. 5,298,071 and 6,197,423 as discussed above, physical deformation of the fiber surface was also believed to improve the fiber-concrete bond. U.S. Pat. No. 4,297,414 of Matsumoto, as another example, taught the use of protrusions and ridges to enhance bond strength. Other surface treatments, such as the use of embossing wheels to impose patterns on the fiber, were also used for improving fiber-concrete bond. Fiber designers have even bent fibers into sinusoidal wave shapes to increase the ability of fibers to resist being pulled out from concrete. However, the present inventors realized that increased structural deformations in the fiber structure may actually enhance opportunities for unwanted fiber balling to occur.\nAgainst this background, the present inventors see a need for novel polymeric synthetic reinforcing fibers having ease of dispersibility in concrete so as to avoid fiber balling and to achieve intended fiber dosage rates, while at the same time to provide strength and fracture toughness in matrix materials and particularly brittle materials such as concrete, mortar, shotcrete, gypsum fireproofing, and the like."}
{"text": "Disposable absorbent articles such as disposable diapers, incontinence pads, training pants, and catamenial napkins generally include an absorbent core for receiving and holding bodily exudates. The absorbent core typically includes a fibrous web, which can be a nonwoven, airlaid web of natural or synthetic fibers, or combinations thereof. Fibrous webs used in such absorbent articles also often include certain absorbent gelling materials usually referred to as “hydrogels”, “superabsorbent” or “hydrocolloid” materials to store large quantities of the discharged bodily liquids. These materials absorb through capillary or osmotic forces, or a combination of both.\nMany absorbent articles such as catamenial pads, adult incontinent products, and diapers are subject to leakage due to an inability to absorb second and subsequent discharges of liquid even if the first liquid discharge has been effectively absorbed. Leakage due to second and subsequent discharges is especially prevalent during the night, when users commonly experience multiple discharges before the absorbent article is changed. One reason for the inability of many absorbent articles to adequately handle multiple discharges of liquid is the inability of the absorbent core to transport discharged liquid away from the region of discharge once the absorbent capacity of that region has been reached. Thus, the overall performance of the absorbent article is limited by the inability to transport the liquid to the farthest reaches of the absorbent core.\nAn alternative absorbent material capable of providing capillary liquid transport is open-celled polymeric foam. Appropriately made open-celled polymeric foams provide features of capillary liquid acquisition, transport, and storage required for use in high performance absorbent cores for absorbent articles such as diapers. Shaped or contoured absorbent cores made from such open-celled foam materials having particularly desirable liquid transport characteristics are disclosed in U.S. Pat. No. 5,147,345 issued to Young et al. on Sep. 15, 1992. The Young et al. '345 absorbent core includes both a liquid acquisition/distribution member and a liquid storage/redistribution member. The liquid acquisition/distribution member is positioned within the absorbent article in such a way as to receive or contact aqueous bodily liquid which has been discharged into the absorbent article by the wearer of the article. The liquid storage/redistribution member, in turn, is positioned within the article to be in capillary liquid communication with the liquid acquisition/distribution member.\nAbsorbent cores providing the desirable absorbent characteristics of the Young et al. '345 patent in an exemplary multi-piece configuration are disclosed in U.S. Pat. No. 5,906,602 issued to Weber et al. on 25 May 1999, which describes shaped absorbent cores having a front panel and a back panel in capillary liquid communication with a center section. The center section includes material generally suitable for liquid acquisition/distribution, while the front and back panels include material generally suitable for liquid storage/redistribution.\nDespite the advances in absorbent articles and in liquid handling absorbent core materials, absorbent articles having multiple absorbent core components as well as those having unitary absorbent cores are generally designed for single use wear. Once the storage/redistribution member is saturated with bodily discharges, such as urine, the entire absorbent article is generally discarded and replaced. Often parts of the absorbent article are still usable, and except for being unitary with the absorbent core, these parts could be used further. In addition to the added cost and waste associated with discarding reusable materials, it is often inconvenient to remove and replace the entire absorbent article when absorbent core components are saturated.\nAbsorbent articles having removable absorbent inserts and thereby being potentially usable for more than a single use are known in the art. For example, U.S. Pat. No. 4,597,761 to Buell, issued Jul. 1, 1986, discloses a disposable absorbent insert for use inside an over-garment such as a conventional reusable diaper, or a disposable diaper. Once the absorbent insert becomes saturated it may be removed and discarded. The absorbent article may then be reused with a fresh absorbent insert. However, because the absorbent insert is removable only from the body side of the article, the absorbent article must be removed from the wearer in order to remove the insert. Therefore, the removal of the absorbent insert is often inconvenient and time consuming.\nAccordingly, it would be desirable to provide an absorbent article having a replaceable absorbent core component wherein the absorbent core component can be replaced without having to remove the absorbent article from the wearer.\nAdditionally, it would be desirable to provide an absorbent article having a replaceable absorbent core component and an apertured garment-facing layer, allowing a saturated component of the absorbent core to be removed through the aperture, thereby exposing an unsaturated component and allowing for prolonged use of the reusable portions of the absorbent article.\nFurthermore, it would be desirable to provide an absorbent article having a replaceable absorbent core component disposed between a wearer-facing layer and a garment-facing layer, configured such that access to the replaceable absorbent core component is gained by separating the wearer-facing layer and the garment-facing layer in a predetermined area to form an opening."}
{"text": "The rapid decrease in viscosity of liquids with increasing temperature is well-known. Ideally, for many applications (automobile lubricants, etc.) it would be desirable to solve this problem so that liquid viscosity would be insensitive to temperature. Alternatively, it might be desirable to provide liquid systems whose viscosities actually increase with temperature. It is true that with selected polymeric additives, it has been possible to reduce substantially the viscosity change with temperature which does occur with most oils and similar systems. These polymer additives, known as viscosity index improvers (or V.I. Improvers), are generally high molecular weight polymers.\nThe way in which these additives function can be summarized very briefly. In effect, they perform two functions, i.e., thickening, which merely increases fluid viscosity; and Viscosity Index (V.I.) improvement, which corresponds to limited thickening at ambient temperatures and a corresponding greater thickening at elevated temperatures. This can be accomplished by utilizing a polymeric additive which is poorly solvated by the liquid at ambient temperatures; however, at elevated temperatures, the polymer is more highly solvated such that the polymer expands and is a relatively more effective thickener.\nWhile these V.I. Improvers have proven successful commercially, it is important to note that their effect at reducing viscosity changes with temperatures is rather mild. For a typical base oil containing a suitable V.I. Improver, the kinematic viscosity will still decrease by a factor of from 5 to 10 as the temperature increases from 30.degree. to 100.degree. C. Obviously, if it is desired to hold the viscosity roughly constant with such temperature changes, current technology has not offered an appropriate additive system.\nU.S. Pat. No. 3,396,136 describes how copolymers of alkenylaromatic sulfonic acids, when properly neutralized, can be employed as thickeners for nonpolar solvents. Those metal sulfonate systems have been shown to be very effective; however, when employed as 2 component systems (i.e., ionic polymer plus nonpolar solvent), the variation of viscosity with increased temperature is very conventional and predictable. That is, the solution viscosity decreases markedly as temperature is increased.\nU.S. Pat. No. 3,396,136 further teaches \"in situ\" neutralization of the sulfonic acid polymer which, under some conditions, can result in the availability of a small amount of polar cosolvent, i.e., a solvent for the sulfonate groups about equal in amount to the amount of sulfonate groups which are present. This amount of polar cosolvent is not within the critical limits of the instant invention, which require amounts of the third component (which interacts with the ionomeric groups of the ionomer copolymer) at levels which range from 10 to 600 times the molar equivalence of ionic groups. This level of cosolvent is about one to two orders of magnitude or more higher than employed in the cited art. This amount of such cosolvent is required to obtain the unusual and unexpected viscosity-temperature behavior which is observed. In addition, the cited patent is restricted to aromatic sulfonate polymers. The instant invention describes other polymers such as co- or terpolymers of acrylates or methacrylates with a sulfonate-containing monomer.\nU.S. Pat. No. 3,666,430 teaches the gelling of organic liquids by the interaction of polar \"associative bonds\" which includes hydrogen bonding and \"ionic crosslinking.\" Again this patent specifies that two components are necessary--the associating polymer (or polymers in some cases) and the nonpolar organic liquid. There is no mention of a third polar cosolvent except to point out that such polar liquids should not be present. Specifically, this patent states (column 2, line 7) that the hydrocarbon liquids to which this invention is to be applied should not contain a substantial portion of a miscible protolytic liquid such as methanol. It is clear that the language of this patent limits this invention to gels and further, that any amount of polar liquids which are present to an extent where they disrupt those gels are undesirable. The instant invention is distinct from that cited in that amounts of such polar compounds, as will break up gel at ambient conditions are required, and in fact, the most preferred state is free of any said gel at ambient temperatures.\nU.S. Pat. No. 3,679,382 teaches the thickening of aliphatic hydrocarbons with synthetic organic polymers which contain olefinically unsaturated copolymerizable acids, amides, hydroxyacrylic esters, sulfonic acids, etc. It is emphasized in this patent (column 3, line 72) that it is critical in the preparation of such polymers that no surface active agent, catalyst or other additive be employed which introduces a metallic ion into the system. Therefore, it is preferred to employ ammonium or amine salts. It is clear that this invention (U.S. Pat. No. 3,679,382) specifically precludes the use of metallic counterions and is directed towards amine or ammonium derivatives. Only metallic counterions are effective in the instant invention and, in fact, attempts to employ amine derivatives have not resulted in the products which are the objective of this invention. Finally, this cited patent does describe (column 7, lines 13-19) that the addition of alcohols will reduce the viscosity of the thickened hydrocarbon and alter flow characteristics thereof.\nOne class of oil soluble polymers known to possess good low temperature viscometrics and to be practically of great interest in the V.I. additive area is that of the acrylate or methacrylate polymers. These long chain acrylate systems (alkyl groups of about 10 to 12 carbons or so) are sold commercially and have been very successful in achieving a suitable balance of thickening behavior and a combination of good low temperature fluidity and high temperature viscosity at reasonable polymer levels. This invention proposes a class of materials that will have comparable or superior viscometrics as compared to the acrylate or methacrylate family of materials. This preferred class of polymers is based on combining sulfonate containing monomers with acrylates and methacrylates of suitable chain length by appropriate polymerization techniques.\nThe instant invention represents an improvement over the process as described in U.S. Pat. No. 3,931,021, wherein the class of co- or terpolymers, as described in the instant invention, show improved low temperature viscometric properties. The co- and terpolymers of the instant invention are restricted to a class of polymers comprising copolymers and terpolymers derived from a metal sulfonate monomer and an alkyl acrylate or alkyl methacrylate, wherein the alkyl group has at least eight carbon atoms."}
{"text": "1. Field of the Invention\nThis invention relates to a superconducting energy stabilizer having multiple load connections employing DC-DC conversion. In particular, the invention relates to an energy stabilizer which utilizes energy stored in a superconducting inductive energy storage device to stabilize the delivery of energy to a DC or AC load or to retrieve and store regenerated energy supplied by the load. Additionally, the invention can supply stored energy to a utility or to an industrial electrical DC or AC distribution system to maintain the level of energy supplied to the consumers.\n2. Background of the Prior Art\nPrior art superconducting energy storage systems have been directed to the storage of energy derived from an alternating current power system and to the return of the stored energy to the alternating current (AC) power system. Typically, energy is stored during periods of low energy demand when surplus energy exists and is returned to the power system when energy demand is high.\nOne such system appears in Peterson et al., U.S. Pat. No. 4,122,512. Peterson et. al. uses a reversible alternating to direct current converter to convert the alternating current obtained from the AC power system into direct current for storage in a superconducting magnet and for converting the stored direct current into alternating current for return to the AC power system. This type of system, while providing for the efficient storage and return of energy to the AC power system, cannot perform the simultaneous storage and release of energy.\nAnother superconducting energy storage system described in Higashino, U.S. Pat. No. 4,695,932, discloses an energy storage circuit which converts three-phase alternating current to direct current for storage in a superconducting energy storage coil. A DC capacitor and chopper circuit are used to control the amount of direct current in the superconducting energy storage coil. According to Higashino, this configuration allows the current capacity of the converter to be scaled down in accordance with service power established by the current rating of the coil, and also allows a reduction of operation losses.\nKnown superconducting energy storage systems provide for the efficient storage and return of energy but cannot perform the simultaneous storage and release of energy. In addition, Higashino has a practical problem of charging the superconductive energy storage coil with a capacitor whose voltage must correspond to the voltage of the AC power system. For a 480 Volt AC system this would be around 650 volts DC. This voltage is too high to charge the superconducting energy storage coil near its maximum rated current, unless the coil is significantly oversized to accommodate the higher voltage value. Normally, a voltage of ten to twenty volts would be more appropriate for charging the coil. Even if Higashino's system used a step down transformer at the AC line in order to drop the capacitor bank voltage, the capacitor bank would need to provide a higher current which results in higher losses and costs due to a bigger magnet current and size.\nIn addition, these systems do not provide a means to boost the current delivered by the superconducting energy storage coil. If a utility or industrial electrical distribution system needs higher current than the magnet current, the known system using a chopper and a capacitor can not provide such a function.\nVisser et al., U.S. Pat. No. 4,962,354, issued to Superconductivity, Inc. on Oct. 9, 1990 describes a superconducting voltage stabilizer. The superconducting voltage stabilizer comprises an AC/DC converter, a voltage regulator, an energy storage cell and a superconducting energy storage coil. Alternating current is converted to direct current and is stored in a superconducting energy storage coil. The stored direct current is released to an energy storage cell which in turn delivers the energy to a load. Energy is released in a controlled fashion according to the energy requirements of the load to thereby stabilize the load's operation. This patent, U.S. Pat. No. 4,962,354, is incorporated herein by reference."}
{"text": "The present invention relates to a support pad for a long hole drill for metal machining with improved properties specially for STS (Single Tube System) system and for drilling tools based on the Ejector® system. The invention also relates to a tool head equipped with a support pad, as well as a method of fabricating such a support pad.\nDrilling in metals is divided generally in two types: long hole drilling and short hole drilling. By short hole drilling is meant generally drilling to a depth of up to 3-5 times the drill diameter.\nIn short hole drilling, the demands are not great, enabling the use of simple helix drills formed either of solid cemented carbide or as solid tool steel or of tool steel provided with cemented carbide inserts.\nLong hole drilling, however, puts large demands on good chip formation, lubrication, cooling and chip transport. This is achieved through specially developed drilling systems with specially designed drilling heads fastened to a drill rod and fulfilling the above mentioned demands. The drill head may be of solid cemented carbide but is generally of tool steel provided with a number of cutting inserts of cemented carbide placed in such a way that they together form the necessary cutting edge. In the center of the head, a tough grade of insert is generally used and on the periphery a more wear resistant one. The cutting inserts are brazed or mechanically clamped. Generally the cutting inserts are brazed on small-diameter drills and mechanically secured on larger-diameter drills. In addition, the head is provided with support pads.\nCutting inserts and support pads are made of cemented carbide with various contents of WC, Co and cubic carbides TiC, TaC, NbC depending on the application. Cutting inserts as well support pads are in addition coated with one or more wear resistant layers. Generally, the support pads are coated with a yellow layer of TiN. A long hole drill head with a geometrical configuration of inserts and support pads as outlined above is disclosed in U.S. Pat. No. 6,682,275 the disclosure of which is hereby incorporated by reference.\nIt is a primary purpose of the invention to provide a support pad and a long hole drill head in combination with such a support pad wherein the useful life and strength is substantially increased, especially when drilling in a titanium-based workpiece.\nA second object of the invention is to provide an improved support head and long hole drill head having improved properties for the machining of heat resistant material, titanium based material and cast iron, as well as in the machining of nickel based alloys."}
{"text": "This application claims the benefit of priority under 35 U.S.C. xc2xa7119 to Japanese Patent Application No. 2001-74236 filed on Mar. 15, 2001, the entire contents of which are incorporated by reference herein.\n1. Field of the Invention\nThe present invention relates to a semiconductor memory device which dynamically stores data with using a channel body of a transistor as a storage node.\n2. Related Background Art\nIn a related DRAM, a memory cell is composed of an MOS transistor and a capacitor. The scale-down of the DRAM has been remarkably advanced by the adoption of a trench capacitor structure and a stacked capacitor structure. At present, the cell size of a unit memory cell is scaled down to an area of 2Fxc3x974F=8F2, where F is a minimum feature size. Namely, the minimum feature size F decreases with the advance of generation, and when the cell size is generally taken to be xcex1F2, xcex1 coefficient a also decreases with the advance of generation. Thus, at the present of F=0.18 xcexcm, xcex1=8 is realized.\nIn order to hereafter secure the trend of cell size or chip size which is the same as before, it is demanded to satisfy xcex1 less than 8 in F less than 0.18 xcexcm and further satisfy xcex1 less than 6 in F less than 0.13 xcexcm, and together with microfabrication, the formation of cell size of the possible small area becomes a large problem. Accordingly, various proposals for decreasing the size of the one memory cell with the one transistor and one capacitor to 6F2 or 4F2 are made. However, practical use is not easy since there are a technical difficulty that the transistor has to be a vertical type, a problem that electric interference between adjacent memory cells increases, and in addition difficulties in terms of manufacturing technology including fabrication, film formation, and the like.\nIt is therefore an object of the present invention to provide a semiconductor memory device which includes memory cells with simple transistor structure and which can store data dynamically.\nAccording to one aspect of the present invention, a semiconductor memory device including MIS transistors to constitute memory cells, each of the MIS transistors comprising:\na semiconductor layer;\na source region formed in the semiconductor layer;\na drain region formed apart from the source region in the semiconductor layer, the semiconductor layer between the source region and the drain region serving as a channel body in a floating state;\na first gate which forms a channel in the channel body;\na second gate formed so as to control a potential of the channel body by a capacitive coupling; and\na high concentration region formed in the channel body on the second gate side, impurity concentration of the high concentration region being higher than that of the channel body,\nwherein the MIS transistor has a first data state in which the channel body has a first potential and a second data state in which the channel body has a second potential.\nAccording to another aspect of the present invention, a semiconductor memory device including MIS transistors to constitute memory cells, each of the MIS transistors comprising:\na semiconductor layer;\na source region formed in the semiconductor layer;\na drain region formed apart from the source region in the semiconductor layer, the semiconductor layer between the source region and the drain region serving as a channel body in a floating state; and\na first gate which forms a channel in the channel body,\nwherein characteristics of the MIS transistor in the case where a channel current flows from the source region to the drain region is different from characteristics of the MIS transistor in the case where a channel current flows from the drain region to the source region, even when the same potential is applied to the first gate, and\nwherein the MIS transistor has a first data state in which the channel body has a first potential and a second data state in which the channel body has a second potential, the first data state being set by the impact ionization generated near a drain junction or by a drain leakage current caused by the first gate, the second data state being set by sending a forward bias current between the drain region and the channel body.\nAccording to a further aspect of the present invention, a semiconductor memory device including MIS transistors to constitute memory cells, each of the MIS transistors comprising:\na semiconductor layer;\na source region formed in the semiconductor layer;\na drain region formed apart from the source region in the semiconductor layer, the semiconductor layer between the source region and the drain region serving as a channel body in a floating state; and\na gate which forms a channel in the channel body,\nwherein the MIS transistor has a first data state in which the channel body has a first potential and a second data state in which the channel body has a second potential, the first data state being set by a drain leakage current caused by the gate to which a negative potential is applied, the second data state being set by sending a forward bias current between the drain region and the channel body."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical signal transmission system, and more particularly, to an optical signal transmission system in analog data service such as a bidirectional cable television (CATV) in which high carrier to noise ratio (CNR) characteristic is requested.\n2. Description of the Related Art\nIn the next generation CATV systems discussed recently from various viewpoints, a digital video signal needs to be transmitted for broad band bidirectional communication services by video on demand (VOD) signals from subscribers services, as well as conventional broadcasting image signals (vestigial sideband--analog modulation (VSB-AM) signals). Further, it is importance to implement such a transmission system with low cost.\nAs a transmission system which satisfy the needs, there is a subcarrier multiplexed (SCM) system which is described by Robert Olshansky et al. in \"Subcarrier Multiplexed Lightwave Systems for Broad-Band Distribution\", (Journal of Lightwave Technology, Vol. 7, No. 9, September 1989, pp. 1329-1341) as reference 1 and by Winston I. Way in \"Subcarrier Multiplexed Lightwave System Design Considerrations For Subscriber Loop applications\", (Journal of Lightwave Technology, Vol. 7, No. 11, November 1989, pp. 1806-1818) as reference 2.\nThis SCM system is an optical signal transmission system using a single mode optical fiber as a transmission medium. Base band analog or digital signals for a plurality of channels, e.g., 60 channels are first high-frequency converted or up-converted using local oscillation circuits (LOs) having different frequencies and then frequency-division-multiplexed. These up-converted signals are synthesized or combined and a laser unit is modulated in accordance with the synthesized high-frequency signal or a combined signal. The frequency of an oscillation signal from the local oscillation circuit is a subcarrier for an optical carrier frequency. In the receiving unit, one of the frequency-division-multiplexing (FDM) channels is selected by a subscriber through a tuning operation of a local oscillation circuit, like the usual tuning operation in a television receiver and radio. A high frequency/microwave signal of the selected FDM channel is low-frequency-converted or down-converted to restore an original base band analog or digital signal.\nIn the above SCM system, in a case where analog image signal for 60 channels and a VOD signal as a digital video signal are simultaneously transmitted, there is a problem in that noise is generated due to clipping distortion so that a bit error rate of the VOD signal is increased because an effective optical modulation exceeds 20% for the image signals for 60 channels.\nIn order to solve the problem, an optical SCM transmission system is proposed in the proceedings of the 1994 spring Conference of IEICE (B-1128) (reference 3). In the reference 3, there is discussed an optimal method in which the bit error rate due to the clipping distortion can be reduced in the transmission of the combined signal composed of a VSB-AM signal as analog image signals and multivalue digital carrier signal (16QAM) as the digital video signal (to be referred to as a VOD signal) and it is concluded that it is effective to separate an image signal transmission system and a VOD signal transmission system from each other such that the VSB-AM signal and the 16QAM signal are transmitted in parallel.\nFIG. 1 is a block diagram showing a conventional optical signal transmission system described in the reference 3. Referring to FIG. 1, the conventional optical signal transmission system includes an optical signal transmitting unit 101 for outputting an optical signal LF60 corresponding to image signals 60VA for 60 channels and an optical signal LV corresponding to a VOD signal VD independently, optical fibers 102 and 104 for transferring the optical signals LF60 and LV, respectively, and an optical signal receiving unit 103 for receiving the optical signals LF60 and LV and detecting a reception combined signal.\nThe optical signal transmitting unit 101 includes frequency division multiplexing (FDM) circuits 111 for frequency-division-multiplexing the image signals 60VA for 60 channels to produce an FDM signal F60, a semiconductor laser (LD) 112 which is subjected to analog-modulation based on the FDM signal F60 and outputs an optical signal LF60 on the optical fiber 102 via an optical isolator 113, the optical isolator 113 for preventing reflected light of the optical signal LF60 from being inputted into the semiconductor laser 112, an FDM circuit 114 for frequency-division-multiplexing the VOD signal VD to produce an FDM signal FV, and a semiconductor laser (LD) 115 which is subjected to analog-modulation based on the FDM signal FV and outputs an optical signal LV on the optical fiber 104. The receiving unit 103 includes an optical signal detector (PD) 121 coupled to the optical fiber 102, for receiving and photoelectrically converting the optical signal LF60 into a reception signal R60 and an optical signal detector (PD) 122 coupled to the optical fiber 104, for receiving and photoelectrically converting the optical signal LV into a reception signal RV.\nNext, the operation of the conventional optical signal transmission system will be described with reference to FIGS. 2A and 2B which indicate the examples of analog-modulation of a semiconductor laser (LD). As shown in FIG. 2A, the analog-modulation of the semiconductor laser is a so-called small amplitude modulating operation in which the semiconductor laser is supplied with a current which satisfies the oscillation condition of the semiconductor laser, i.e., a constant bias current higher than a threshold current and a modulation current is superposed on the bias current. If the amplitude of the modulation current becomes greater so that the negative peak of the modulation current is decreased lower than the threshold current as shown in FIG. 2B, a part of the modulated light output signal, the negative peak in the example is cut so that the clipping distortion is caused. Also, when the output of the semiconductor laser is saturated at the positive peak, the clipping distortion is caused. Thus, the clipping distortion is caused when the optical modulation of the semiconductor laser is more than 100%.\nIn this conventional optical signal transmission system, since an FDM signal F60 corresponding to the image signals for 60 channels needs to be transmitted, the optical signal outputted from the single semiconductor laser 112 is modulated in accordance with the FDM signal F60. In this case, an effective optical modulation .mu.e is represented by the following equation for the FDM signal F60. ##EQU1## where N is the number of image signals which constitute the FDM signal F60 and .mu.i is the optical modulation index for the image signal VAi, A relation between the effective modulation index .mu.e and the clipping distortion is experimentally obtained. For instance, when an optical modulation index is 5% for an image signal for one channel and image signals for 60 channels are supplied to the semiconductor laser 112 as the FDM signal F60, the effective optical modulation index .mu.e of the semiconductor laser 112 is about 27%.\nFIG. 3 is a graph showing a relation of a bit error rate (BER) to a carrier to noise ratio (CNR) when image signals of a plurality of channels and a VOD signal of one channel are simultaneously transmitted through modulation by a single semiconductor laser. It can be seen that a bit error is abruptly increased if the effective optical modulation index exceeds 20%. This is because the impulse noise due to the clipping distortion adversely influences to the modulated signal."}
{"text": "Many different polymers and materials have been added to polymer compositions in attempting to enhance the composition's impact strength or maintain the impact strength while enhancing other properties. For example, U.S. Pat. No. 5,118,753 (Hikasa et al.), incorporated herein by reference, discloses thermoplastic elastomer compositions said to have low hardness and excellent flexibility and mechanical properties consisting essentially of a mixture of an oil-extended olefinic copolymer rubber and an olefinic plastic. The olefinic plastic is polypropylene or a copolymer of polypropylene and an .alpha.-olefin of 2 or more carbon atoms. Modern Plastics Encyclopedia/89, mid October 1988 Issue, Volume 65, Number 11, pp. 110-117, the disclosure of which is incorporated herein by reference, also discusses the use of various thermoplastic elastomers (TPEs) useful for impact modification. These include: elastomeric alloys TPEs, engineering TPEs, olefinic TPEs (also known as thermoplastic olefins or TPOs), polyurethane TPEs and styrenic TPEs.\nThermoplastic olefins (TPOs) are often produced from blends of an elastomeric material such as ethylene/propylene rubber (EPM) or ethylene/propylene diene monomer terpolymer (EPDM) and a more rigid material such as isotactic polypropylene. Other materials or components can be added into the formulation depending upon the application, including oil, fillers, and cross-linking agents. TPOs are often characterized by a balance of stiffness (modulus) and low temperature impact, good chemical resistance and broad use temperatures. Because of features such as these, TPOs are used in many applications, including automotive facia and wire and cable operations, rigid packaging, molded articles, instrument panels, and the like.\nUnion Carbide Chemicals and Plastics Inc. announced in 1990 that they have developed a new cost effective class of polyolefins trademarked Flexomer™ Polyolefins that could replace expensive EPM or EPDM rubbers. These new polyolefins are said to have bridged the gap between rubbers and polyethylene, having moduli between the two ranges. Modulus of the rubber and of the formulation is not, however, the only criteria for evaluating a TPO formulation. Low temperature impact performance, sometimes measured by Gardner Impact at −30° C. also is sometimes important to a TPO composition's performance. According to the data contained in FIG. 4 of the paper “Flexomer™ Polyolefins: A Bridge Between Polyethylene and Rubbers” by M. R. Rifi, H. K. Ficker and M. A. Corwin, more of the Flexomer™ Polyolefin needs to be added into the TPO formulation in order to reach the same levels of low temperature Gardner Impact performance as the standard EPM rubber, thus somewhat negating the benefits of the lower cost EPM/EPDM replacement. For example, using the data of FIG. 4 of the Rifi et al paper, about 20% (by weight) of the EPM in polypropylene gives a Gardner Impact of about 22 J. at −30.degree. C., while the same amount of Flexomer™ Polyolefin gives a −30° C. Gardner Impact of about 13 J.\nIn a paper presented on Sep. 24, 1991 at the 1991 Specialty Polyolefins Conference (SPO '91) (pp. 43-55) in Houston, Tex., Michael P. Jeffries (Exxpol Ethylene Polymers Venture Manager of Exxon Chemical Company) also reports that Exxon's Exact™ polymers and Plastomers can be blended into polypropylene for impact modification. Exxon Chemical Company, in the Preprints of Polyolefins VII International Conference, page 45-66, Feb. 24-27 1991, also disclose that the narrow molecular weight distribution (NMWD) resins produced by their EXXPOL™ technology have higher melt viscosity and lower melt strength than conventional Ziegler resins at the same melt index. In another recent publication, Exxon Chemical Company has also taught that NMWD polymers made using a single site catalyst create the potential for melt fracture (“New Specialty Linear Polymers (SLP) For Power Cables,” by Monica Hendewerk and Lawrence Spenadel, presented at IEEE meeting in Dallas, Tex., September, 1991).\nIt is well known that narrow molecular weight distribution linear polymers disadvantageously have low shear sensitivity or low I10/I2 value, which limits the extrudability of such polymers. Additionally, such polymers possessed low melt elasticity, causing problems in melt fabrication such as film forming processes or blow molding processes (e.g., sustaining a bubble in the blown film process, or sag in the blow molding process etc.). Finally, such resins also experienced surface melt fracture properties at relatively low extrusion rates thereby processing unacceptably and causing surface irregularities in the finished product.\nThus, while the development of new lower modulus polymers such as Flexomer™ Polyolefins by Union Carbide or Exact™ polymers by Exxon has aided the TPO marketplace, there continues to be a need for other more advanced, cost-effective polymers for compounding with thermoplastics (e.g., polyolefins such as polypropylene or HDPE) to improve or maintain modulus and/or impact performance at room temperature or below.\nFormulated compositions have now been discovered to have this combination of good low temperature impact performance and modulus. The compositions comprise:\nA) a thermoplastic selected from the group consisting of thermoplastic polyurethanes, polyvinyl chlorides, styrenics, engineering thermoplastics, and polyolefins, and\nB) an impact-modifying amount of at least one ethylene/α-olefin interpolymer."}
{"text": "Broadband communications have become an increasingly prevalent form of electromagnetic information exchange and coaxial cables are common conduits for transmission of broadband communications. Coaxial cables are typically designed so that an electromagnetic field carrying communications signals exists only in the space between inner and outer coaxial conductors of the cables. This allows coaxial cable runs to be installed next to metal objects without the power losses that occur in other transmission lines, and provides protection of the communications signals from external electromagnetic interference. Connectors for coaxial cables are typically connected onto complementary interface ports to electrically integrate coaxial cables to various electronic devices and cable communication equipment. Connection is often made through rotatable operation of an internally threaded nut of the connector about a corresponding externally threaded interface port. Fully tightening the threaded connection of the coaxial cable connector to the interface port helps to ensure a ground connection between the connector and the corresponding interface port. However, often connectors are not properly tightened or otherwise installed to the interface port and proper electrical mating of the connector with the interface port does not occur. Moreover, structure of common connectors may permit loss of ground and discontinuity of the electromagnetic shielding that is intended to be extended from the cable, through the connector, and to the corresponding coaxial cable interface port. Hence a need exists for an improved connector for ensuring ground continuity between the coaxial cable, the connector structure, and the coaxial cable connector interface port."}
{"text": "The invention relates to defined tetrakisazo dyes and mixtures thereof with defined trisazo dyes, their production and their use, especially for the dyeing of leather.\nIn U.S. Pat. No. 4,479,906 there are described trisazo dyes of the formula \nin which\nAO is the radical of an aromatic-carbocyclic diazo component,\nBO is the radical of a coupling component which contains hydroxyl and/or amino groups,\nn is the integer 1 or 2,\none XO is hydroxyl and the other XO is amino,\nZO is \nTO is hydrogen, methyl, ethyl, benzyl, xcex2-hydroxyethyl, xcex2-hydroxypropyl, xcex2-cyanoethyl, xcex2-carboxyethyl or xcex2-carbalkoxyethyl of a total of two to five carbon atoms,\nRO is hydrogen, hydroxyl, chloro, bromo, nitro, methyl, ethyl or carboxyl and\nR1O is hydrogen, methyl, ethyl, chloro, bromo, cyano, nitro, hydroxysulphonyl, carboxyl or alkoxycarbonyl,\nthere being described in particular also some tetrakisazo dyes in which in the above formula n is 2, the two sulpho groups are in the positions meta to each XO (as corresponds to H-acid), ZO is a radical of formula (2). AO is the radical of an aminoazobenzene diazocomponent and BO is the radical of certain aminobenzene coupling components. These dyes are described also for the dyeing of leather in green to black shades.\nIt has now been found that tetrakisazo dyes of the below indicated structure of formula (I) are particularly valuable anionic dyes for obtaining dyeings of intensive dark shades, in particular for the penetration dyeing of leather in dark shades of high fastnesses, especially when combined with the below defined trisazo dyes of formula (X).\nThe invention relates to the defined dyes and dye mixtures, their production and their use.\nThe invention, thus, provides a tetrakisazo dye of formula \nwherein\none of the two symbols X and Y signifies xe2x80x94OH and the other signifies xe2x80x94NH2,\nZ signifies a radical of formula\nxe2x80x94W1xe2x80x94Nxe2x95x90Nxe2x80x94B2 xe2x80x83xe2x80x83(xcex1) \nor, if Y signifies xe2x80x94OH, also a radical of formula\nxe2x80x94W2xe2x80x94Nxe2x95x90Nxe2x80x94D xe2x80x83xe2x80x83(xcex2),\nQ1 signifies xe2x80x94SO2xe2x80x94 or xe2x80x94COxe2x80x94,\nQ2 signifies the direct bond or, if Q1 is xe2x80x94SO2xe2x80x94, also a group xe2x80x94SO2xe2x80x94,\nB1 signifies the radical of a coupling component Hxe2x80x94B1,\nB2 signifies the radical of a coupling component Hxe2x80x94B2,\nD signifies the radical of a diazo component Dxe2x80x94NH2,\nW1 signifies the radical of a bicyclic bisdiazo component,\nW2 signifies the radical of a bicyclic middle component and\nM signifies hydrogen or a cation,\nor a mixture of two or more tetrakisazo dyes of formula (I).\nFormula (I) comprises in particular the following formulae \nThus, the invention provides in particular a dye of formula (Ia), (Ib) or (Ic) or a mixture of two or more thereof.\nA particular further feature of the invention is represented by a mixture of at least one dye of the above formula (I) and at least one trisazo dye of formula \nin which B1, D, Q1, Q2 and M have the significances indicated above for formula (I).\nThe dyes of the invention may be produced by means of diazotization and coupling reactions, in order to form the azo groups, in a manner conventional per se. In particular, a process for the production of the tetrakisazo dyes of the invention is characterized in that for the production of a dye of formula (Ia) the bisdiazocompound of a diamine of formula \nis coupled at one side to the 3-position of a bis-coupling component of formula \nthe obtained azodiazocompound is either first coupled to a coupling component Hxe2x80x94B1, to give an intermediate of formula \nand the diazocompound of an amine of formula\nDxe2x80x94Nxe2x95x90Nxe2x80x94W2xe2x80x94NH2 xe2x80x83xe2x80x83(V) \nis coupled to the 6-position of the intermediate of formula (IV), or the azodiazocompound is first coupled in the 6-position with the diazocompound of an amine of formula (V) and then coupled to a coupling component Hxe2x80x94B1, for the production of a dye of formula (Ib) the bisdiazocompound of a diamine of formula (II) is coupled at one side to the 3-position of a bis-coupling component of formula (III) and the obtained azodiazocompound is either first coupled to a coupling component Hxe2x80x94B1, to give an intermediate of formula (IV), and then the bisdiazocompound of a diamine of formula\nH2Nxe2x80x94W1xe2x80x94NH2 xe2x80x83xe2x80x83(VI) \nis coupled at one side to the 6-position of the intermediate of formula (IV) and the obtained trisazodiazocompound is coupled to a coupling component Hxe2x80x94B2, or the azodiazocompound is first coupled in the 6-position with the bisdiazocompound of a diamine of formula (VI) and then the obtained disazobisdiazocompound is coupled to coupling components Hxe2x80x94B1, and Hxe2x80x94B2 or for the production of a dye of formula (Ic) the bisdiazocompound of a diamine of formula (VI) is coupled at one side to the 3-position of a bis-coupling component of formula (III) and the obtained azodiazocompound is either first coupled to a coupling component Hxe2x80x94B2 to give an intermediate of formula \nand then the bisdiazocompound of a diamine of formula (II) is coupled at one side to the 6-position of the intermediate of formula (VII) and the obtained trisazodiazocompound is coupled to a coupling component Hxe2x80x94B1, or it is first coupled in the 6-position with the bisdiazocompound of a diamine of formula (II) and then the obtained disazobisdiazocompound is coupled to Hxe2x80x94B1, and Hxe2x80x94B2.\nThe compounds of formula (V) may be produced by coupling the diazocompound of a diazocomponent of formula\nDxe2x80x94NH2 xe2x80x83xe2x80x83(VIII) \nto a middle component of formula (IX)\nHxe2x80x94W2xe2x80x94NH2 xe2x80x83xe2x80x83(IX). \nThe trisazo dyes of formula (X) may be synthetized by coupling the diazocompound of a diazocomponent of formula (VIII) to a compound of formula (IV) or to its precursor azodiazocompound, before coupling of the diazoazocompound to Hxe2x80x94B1. If the starting compounds are the same as the corresponding ones employed for the production of the dyes of formula (I), this reaction sequence may be carried out together with the production of the tetrakisazo dye or dye mixture, choosing the molar ratios in a suitable way.\nIn the mixtures of tetrakisazo dyes of formula (I) and trisazodyes of formula (X) the molar ratio of tetrakisazo dye or tetrakisazo dye mixture to trisazo dye or trisazo dye mixture is e.g. in the range of 95/5 to 20/80, prferably 80/20 to 30/70, more preferably 75/25 to 40/60.\nAs diazocomponents of formula (VIII) there may be employed any known diazotizable amines, in particular of the benzene of naphthalene series, preferably containing at least one solubilizing substituent, which preferably is a sulpho group, a carboxy group, a carbamoyl group and/or a sulphamoyl group.\nAs diazo components of formula (VIII) come, in particular, into consideration those of the following formulae \nwherein\nR1 signifies hydrogen, nitril, trifluoromethyl, nitro, xe2x80x94SO3M, xe2x80x94SO2NR4R5, xe2x80x94COOM or xe2x80x94CONR4R5,\nR2 signifies hydrogen, nitro, xe2x80x94SO3M, xe2x80x94SO2NR4R5, trifluoromethyl, nitril, xe2x80x94COOM, xe2x80x94CONR4R5, C1-4-alkyl, C1-4-alkoxy, halogen or C1-2-mercaptoalkyl,\nR3 signifies hydrogen, C1-4-alkyl, C1-4-alkoxy, halogen, C1-2-mercaptoalkyl, xe2x80x94NHxe2x80x94Ac, xe2x80x94NHxe2x80x94COxe2x80x94Oxe2x80x94CH3 or a radical of formula \nR4 signifies hydrogen, C1-4-alkyl, C1-3-alkylene-Rxe2x80x2 or C2-3-3-hydroxyalkyl,\nR5 signifies hydrogen, C1-4-alkyl, C1-3-alkylene-Rxe2x80x2, C2-3-hydroxyalkyl, benzyl or a radical of formula \nRxe2x80x2 signifies nitril, carbamoyl or xe2x80x94COOM,\nR6 signifies hydrogen, nitro or xe2x80x94SO3M,\nR7 signifies hydrogen, methyl, chlorine, nitro, xe2x80x94COOM or xe2x80x94SO3M,\nR8 signifies hydrogen, halogen, nitro, xe2x80x94SO3M, xe2x80x94SO2NR4R5, methylsulphonyl, C1-4-alkyl or xe2x80x94NHxe2x80x94Ac,\nR9 signifies hydrogen, halogen, nitro, xe2x80x94SO3M, xe2x80x94SO2NR4R5, methylsulphonyl, C1-4-alkyl or xe2x80x94NHxe2x80x94Ac,\nR10 signifies nitro, xe2x80x94SO3M or xe2x80x94SO2NR4R5,\nAc signifies an aliphatic acyl group,\np signifies 0 or 1 and\nr signifies 0 or 1,\nR1 preferably signifies hydrogen or xe2x80x94COOM.\nOf the alkyl and alkoxy groups with 1 to 4 carbon atoms, the lower molecular ones are preferred, in particular ethyl, methyl, ethoxy and methoxy. Halogen may in particular be fluorine, chlorine or bromine of which fluorine and, before all, chlorine are preferred.\nWhere in formula (a1) R3 signifies a radical of formula (axe2x80x2) this is preferably in para-position to the amino group, R1 preferably signifies hydrogen, R2 preferably signifies hydrogen or a sulpho group and, if R2 signifies a sulpho group, this is preferably in meta-position to the azo group and R6 and R7 preferably signify both hydrogen or, if R2 signifies hydrogen, R6 signifies preferably a nitro group and R7 signifies preferably a carboxy group or a sulpho group, the two substituents R6 and R7 being placed in the para-position to the imino group and in one of the two ortho-positions to the imino group.\nWhere the symbol R3 does not signify a radical (axe2x80x2), R3 preferably signifies hydrogen and R1 preferably signifies hydrogen or xe2x80x94COOM. According to a preferred feature, in this case in formula (a1) one or both of the positions ortho to the primary amino group and, more preferably, also one or both of the positions meta to the primary amino group are unsubstituted. According to a further preferred feature, if R1 signifies xe2x80x94COOM, this is in position ortho or para to the diazotizable primary amino group.\nR4 preferably signifies hydrogen.\nR5 advantageously signifies hydrogen, methyl, ethyl, hydroxyethyl or a radical of formula (axe2x80x3), in which R7 preferably signifies hydrogen or carboxy.\nThe aliphatic acyl group Ac advantageously signifies the radical of a low molecular aliphatic carboxylic acid, preferably of an alkanoic acid with 2 to 4 carbon atoms, more preferably acetyl or propionyl, of which acetyl is preferred.\nMost preferably the diazo component of formula (a1) is an aminobenzoic acid, in particular anthranilic or p-aminobenzoic acid.\nIn formula (a2 ) preferably at least one of R8 and R9 has a significance other than hydrogen, more preferably R9 has a significance other than hydrogen and R8 signifies hydrogen, a nitro group or a sulpho group.\nThe primary amino group in formula (a4) may be in any of the positions xcex1 and xcex2 of the naphthalene ring and, if p signifies 1, the sulpho group may be in any of the other available positions, preferably so that at least one vicinal position to the amino group is unsubstituted; e.g. if the amino group is in position 1, the sulpho group is preferably in any of the positions 3 to 8, more preferably 4 to 8, and if the amino group is in the position 2, the sulpho group is e.g. in position 1 or in any of the positions 4 to 8, more preferably 5, 6 r 7. A preferred diazocomponent of formula (a4) is naphthionic acid.\nThe compounds (a5) are indicated in the free amino form, as diazo components they are however usually employed directly in the form of the respective commercially available diazonium compounds 1-diazonium-2-naphthol-4-sulphonic acid (in which r=0) and 1-diazonium-6-nitro-2-naphthol-4-sulphonic acid (in which r=1).\nAmong the above diazo components are preferred those, which do not contain a hydroxy group in ortho position to the diazotizable amino group, in particular the compounds of formula (a1) and (a4).\nPreferably Dxe2x80x94NH2 contains at least one hydrophilic substituent, preferably a substituent selected from the group consisting of xe2x80x94SO3M, xe2x80x94SO2NR4R5, xe2x80x94COOM and xe2x80x94CONR4R5, more preferably not more than one hydrophilic substituent per homocyclic aromatic nucleus.\nAs a middle component there is understood a compound that acts as a coupling component and that contains a diazotizable primary amino group or a group convertible to a diazotizable primary amino group and which compound, upon having been coupled with a diazo compound, if required after having converted a substituent convertible to a primary amino group to this primary amino group, is diazotized and the obtained azodiazocompound is then coupled to a further coupling component. As middle components, from which W1 derives, come preferably into consideration middle components containing a primary amino group.\nThe radical W2 and respectively the middle component of formula (IX), from which the radical W2 derives, contains two rings of aromatic character, of which one is a benzene ring and the other may be an optionally condensed second benzene ring (to give a naphthalene radical) or a benzene ring linked over a bridging group or a heterocyclic ring linked to the benzene ring, optionally over a heteroatomic bridging group. The diazotizable primary amino group is suitably a substituent of a benzene ring and the coupling position is suitably activated by at least one electron donor substituent, in particular by an amino group or preferably a hydroxy group. The middle component of formula (IX) preferably is a compound of the aminonaphthalene series (more preferably of the aminonaphthole series), of the 1-(aminophenyl)-5-pyrazolone or -5-aminopyrazol series or of the sulphanilylamino-phenol or -pyrimidone series. Where the middle component is of the pyrazole series, it is preferably a pyrazolone in which the 3-position may be substituted with carboxy or C1-4-alkyl, preferably methyl, and where the primary amino group at the 1-positioned phenyl radical is in meta- or para-position to the pyrazole radical. Where the middle component is of the sulphanilylaminophenol series, the phenolic hydroxy group is preferably in meta-position and the two benzene rings preferably do not contain any further substituents. Where the middle component is of the sulphanilylaminopyrimidine series, this is e.g. a middle component as described in GB-A 2 276 174, among which 6-amino-2-sulphanilylamino-3-H-pyrimidin-4-one and 2-sulphanilylamino-3-H-pyrimidine-4,6-dione are preferred. Among the aminonaphthols are preferred those containing two or preferably only one sulpho group, more particularly H-acid, J-acid and xcex3-acid (the latter being particularly preferred).\nB1 or B2 may be the radical of any coupling component Hxe2x80x94B1 or Hxe2x80x94B2 which may, in particular, be of the benzene, naphthalene, heterocyclic or open-chain methylene-active series, and suitably contains at least one substituent that activates the compound Hxe2x80x94B1 or Hxe2x80x94B2 or coupling, in particular an aromatically bound or enolic, optionally etherified hydroxy group or an optionally substituted amino group, so that the coupling reaction may take place in the corresponding activated position of the molecule Hxe2x80x94B1 or Hxe2x80x94B2.\nSuitable coupling components Hxe2x80x94B1 or Hxe2x80x94B2 are, in particular, those of the following formulae \nwherein\nR13 signifies hydrogen, xe2x80x94OR17 or NHR17,\nR14 signifies xe2x80x94OR17 or xe2x80x94NHR17,\nR15 signifies hydrogen, xe2x80x94SO3M, xe2x80x94SO2NR4R5, xe2x80x94COOM or xe2x80x94CONR4R5.\nR16 signifies hydrogen, xe2x80x94SO3M, xe2x80x94SO2NR4R5, xe2x80x94COOM or xe2x80x94CONR4R5,\nR17 signifies hydrogen, C1-4-hydrogen, C1-4-alkyl, Acxe2x80x2 or a radical of formula \nAcxe2x80x2 signifies the acyl radical of an aliphatic carboxylic acid,\nQ3 signifies xe2x80x94COxe2x80x94, xe2x80x94SO2xe2x80x94 or the direct bond,\nR18 signifies hydrogen, methyl, xe2x80x94NHxe2x80x94Ac, xe2x80x94COOM or xe2x80x94NO2 or, if in formula (bxe2x80x2) Q3 signifies xe2x80x94COxe2x80x94 or xe2x80x94SO2xe2x80x94, also xe2x80x94NH2,\nR19 signifies xe2x80x94OH or xe2x80x94NH2,\nR20 signifies hydrogen, C1-4-alkyl, C1-4-alkoxy, xe2x80x94OH, xe2x80x94NRxe2x80x3Rxe2x80x3xe2x80x2 or xe2x80x94NHxe2x80x94Ac,\nR21 signifies hydrogen, C1-4-alkyl or C1-4-alkoxy,\nRxe2x80x3 and Rxe2x80x2xe2x80x3 independently, signify hydrogen, C1-2-alkyl or C2-3-hydroxy-alkyl,\nR22 signifies hydrogen, sulphonaphthyl or a radical of formula \nG1 signifies hydrogen, halogen, methyl, methoxy or xe2x80x94COOM,\nG2 signifies hydrogen, halogen, trifluoromethyl, nitril, nitro, xe2x80x94COOM, xe2x80x94SO3M or xe2x80x94SO2NR4R5,\nR23 signifies C1-4-alkyl, phenyl, xe2x80x94COOM, xe2x80x94CONR4R5, xe2x80x94COOCH3 or xe2x80x94COOC2H5,\nR24 signifies xe2x95x90O or xe2x95x90NH,\nR25 signifies hydrogen, unsubstituted amino, phenylamino, sulphonaphthyl, open-chain C1-8-alkyl, C6-9-cycloalkyl, carboxy-(C1-4-alkyl), C2-4-alkyl substituted with hydroxy, methoxy, ethoxy or a sulpho group in one of the positions xcex2 to xcfx89, or a radical of formula (bxe2x80x3),\nR26 signifies hydrogen, nitril, acetyl, xe2x80x94COOM, carbamoyl, xe2x80x94SO3M, pyridinio or 2-methyl-pyridinio,\nR27 signifies hydrogen, hydroxy, methyl, carboxy, phenyl, sulphomethyl or carbamoyl,\nR28 signifies hydroxy, primary amino, nitrilamino, thiol or a radical of formula \nR29 signifies hydroxy or primary amino,\nR30 signifies hydroxy or primary amino,\nR31 signifies hydrogen, methyl, chlorine, chloromethyl or chloroacetyl,\nT signifies xe2x80x94Oxe2x80x94, xe2x80x94NHxe2x80x94 or the direct bond,\nR32 signifies naphthyl, sulphonaphthyl, disulphonaphthyl or a radical of formula (bxe2x80x3),\nR33 signifies C1-4-alkyl,\nG3 signifies hydrogen, hydroxy, C1-4-alkoxy, xe2x80x94NHAc, xe2x80x94NH2, di-(C1-4-alkyl)-amino or carboxy-methylamino, and\nG4 signifies hydrogen or methyl,\nand, where in formula (b4) R26 stands for pyridinio or orthomethylpyridinio, any carboxy or sulpho group present in the molecule may be in the form of the anion xe2x80x94SO3xe2x80x94 or xe2x80x94COOxe2x80x94 to form the counter ion in the form of the inner salt.\nThe aliphatic acyl radical Acxe2x80x2 in the significance of R17 may be the radical of any aliphatic carboxylic acid as can be introduced by acylation, in particular of a C2-12-alkanoic primary monocarboxylic acid, preferably such as stated above for Ac, especially C2-4-alkanoyl, most preferably acetyl.\nIf R18 signifies xe2x80x94COOM, Q3 in formula (bxe2x80x2) signifies in particular xe2x80x94COxe2x80x94. If in formula (bxe2x80x2) Q3 signifies the direct bond, R18 preferably signifies hydrogen. If in formula (bxe2x80x2) Q3 signifies xe2x80x94SO2xe2x80x94, R18 preferably signifies methyl, acetylamino or xe2x80x94NH2. If in formula (bxe2x80x2) Q3 signifies xe2x80x94COxe2x80x94, R18 preferably signifies hydrogen or xe2x80x94NO2.\nIn formula (b1) xe2x80x94OR17 preferably signifies hydroxy and xe2x80x94NHR17 preferably signifies xe2x80x94NHR17xe2x80x2, where R17xe2x80x2 signifies hydrogen, methyl, acetyl or a radical of formula (bxe2x80x2). Preferably R14 signifies hydroxy or xe2x80x94NHR17xe2x80x2 and R13 signifies hydrogen or, where R14 signifies xe2x80x94OH, also a group xe2x80x94NHR17xe2x80x2. More preferably either R14 signifies hydroxy and R13 signifies hydrogen or xe2x80x94NHR17 or R14 signifies xe2x80x94NHR17xe2x80x2 and R13 signifies hydrogen, R15 preferably signifies hydrogen, xe2x80x94SO3M, xe2x80x94COOM or xe2x80x94CONH2. R16 preferably signifies hydrogen or xe2x80x94SO3M, more preferably hydrogen. Preferred coupling components of formula (b1) are those in which R13, R15 and R16 signify hydrogen and, in particular, R14 signifies hydroxy. Most preferably (b1) is xcex2-naphthol.\nIf in formula (b2) R20 signifies hydroxy, xe2x80x94NRxe2x80x3Rxe2x80x2xe2x80x3 or xe2x80x94NHxe2x80x94Ac, it is preferably in meta-position to R19 and R21 preferably signifies hydrogen. If R20 signifies C1-4-alkyl or C1-4-alkoxy it may be in any of the available positions ortho, meta and para to R19. More preferably R19 signifies hydroxy. Advantageously R21, signifies hydrogen.\nOf the alkyl and alkoxy radicals with 1 to 4 carbon atoms, also in Hxe2x80x94B1 or Hxe2x80x94B2 the lower molecular ones are preferred (analogously as in Dxe2x80x94NH2), more specifically ethyl, methyl, ethoxy and methoxy. In the C2-3-hydroxyalkyl radicals the hydroxy group is preferably in xcex2position.\nIn formula (b3) R22 preferably signifies a radical of formula (bxe2x80x3). In formula (bxe2x80x3)xe2x80x94in the significance of R22xe2x80x94preferably at least one of G1 and G2 signifies hydrogen, more preferably G1. R23 preferably signifies C1-4-alkyl, more preferably methyl. R24 preferably signifies oxygen.\nThe open-chain C3-8-alkyl radicals in the significance of R25 may be linear or branched, if they contain 6 to 8 carbon atoms they are preferably branched; the cycloalkyl radicals in the significance of R25 are preferably cyclohexyl, which may be substituted with 1 to 3 methyl groups, more preferably it is unsubstituted cyclohexyl. The carboxy-substituted C1-4-alkyl group preferably is carboxymethyl or xcex2-carboxyethyl. The substituent (hydroxy, methoxy, ethoxy, sulpho) at the C2-4-alkyl, in the significance of R25, is preferably in xcex2-position. If R25 signifies a radical of formula (bxe2x80x3) G1 preferably signifies hydrogen and G2 preferably signifies carboxy, sulpho or trifluoromethyl. Preferred significances of R25 are hydrogen, a radical of formula (bxe2x80x3), C1-8-alkyl, C2-3-hydroxyalkyl and C6-9-cycloalkyl.\nR26 preferably signifies hydrogen, a sulpho group or one of the stated nitrogen-containing substituents.\nR27 preferably has a significance other than hydrogen, more preferably methyl.\nIn formula (b5) preferably at least one of R29 and R30 signifies hydroxy, more preferably both R29 and R30 signify hydroxy groups.\nIf in formula (b5) R28 signifies a radical of formula (bxe2x80x3xe2x80x2), xe2x80x94NHxe2x80x94Q3xe2x80x94 preferably signifies a group xe2x80x94NHxe2x80x94SO2xe2x80x94.\nIn formula (b6) the hydroxy group preferably is in position 8. If R31 is other than hydrogen it is preferably in position para to the 8-positioned hydroxy group. R31 preferably signifies hydrogen or methyl, more preferably hydrogen.\nIn formula (b7) T preferably signifies xe2x80x94NHxe2x80x94. More preferably R32 is unsubstituted phenyl and R33 is preferably methyl.\nThe coupling component radical xe2x80x94B1 or xe2x80x94B2 preferably contains up to three cycles (homocyclic rings, heterocyclic rings and optionally a cycloaliphatic ringxe2x80x94a naphthalene radical being calculated as two cycles), more preferably xe2x80x94B1 or xe2x80x94B2 contain each one or two of such cycles.\nThe coupling components Hxe2x80x94B1 and Hxe2x80x94B2 preferably correspond to formula (b1), (b2), (b3) and/or (b7) or are preferably mixtures of (b1) with one of (b2), (b3) or (b7), e.g. in the molar ratio of 1/9 to 9/1, preferably 3/7 to 7/3.\nThe bis-diazo component radicals W1 are bicyclic, in particular they contain two benzenic rings linked to each other over a heteroatomic bridge. They derive from corresponding bis-diazo components containing as a substituent at one of the two rings a primary diazotizable amino group and as a substituent at the other benzenic ring a primary diazotizable amino group or a substituent convertible into such a primary amino group, in particular a nitro group reducible to primary amino or an acylated amino group hydrolyzable to primary amino; more preferably this substituent is, however, primary amino. As bis-diazo components from which W1 derives come, in particular, into consideration those of the formula \nwherein\nT1 signifies a heteroatomic bridging member, in particular xe2x80x94SO2xe2x80x94, xe2x80x94COxe2x80x94, xe2x80x94SO2xe2x80x94NHxe2x80x94, xe2x80x94SO2xe2x80x94NHxe2x80x94SO2xe2x80x94, xe2x80x94COxe2x80x94NHxe2x80x94, xe2x80x94NHxe2x80x94COxe2x80x94NHxe2x80x94, xe2x80x94Sxe2x80x94, xe2x80x94Oxe2x80x94, xe2x80x94NHxe2x80x94, or a methylene bridge.\nRx signifies primary amino, acetylamino or nitro, and\nz signifies 0 or 1.\nIn formula (Aa1) Rx preferably signifies primary amino.\nT1 preferably signifies xe2x80x94SO2NHxe2x80x94, xe2x80x94COxe2x80x94NHxe2x80x94, xe2x80x94NHxe2x80x94 or xe2x80x94SO2NHxe2x80x94SO2xe2x80x94, among which xe2x80x94NHxe2x80x94 and especially xe2x80x94SO2xe2x80x94NHxe2x80x94 are particularly preferred.\nR2 in formula (Aa1) preferably signifies hydrogen.\nWhere T1 signifies xe2x80x94NHxe2x80x94 z preferably signifies 1. Where T1 signifies xe2x80x94COxe2x80x94NHxe2x80x94 or xe2x80x94SO2xe2x80x94NHxe2x80x94 the xe2x80x94NHxe2x80x94 group thereof is preferably linked to the benzenic ring bearing the substituent Rx. Where T1 signifies a bridging group containing a xe2x80x94SO2xe2x80x94 or xe2x80x94COxe2x80x94 group z preferably signifies 0.\nThe substituent Rx, in particular the primary amino group, may be in position meta, or preferably, para to T1.\nThe compounds of formula (VI) preferably correspond to the above formula (Aa1) wherein Rx signifies primary amino, which is more preferably in para position to T1.\nIn formula (II) Q1 preferably signifies xe2x80x94SO2xe2x80x94; Q2 preferably signifies the direct bond.\nAmong the tetrakisazo dyes of the invention are preferred those in which Z signifies a radical of formula (xcex1), especially those of formula (Ib), and particularly preferred are also their mixures with trisazo dyes of formula (X).\nThe diazotization and coupling reactions may be carried out in a manner conventional per se. Diazotization may e.g. be carried out with a nitrite (preferably sodium nitrite) in acidic aqueous medium (preferably in the presence of hydrochloric acid) and at low temperatures, e.g. in the range of xe2x88x925xc2x0 C. to +10xc2x0 C., preferably 0 to 5xc2x0 C. The coupling reaction of the diazonium compounds to the respective coupling components is advantageously carried out at temperatures in the range of xe2x88x925xc2x0 C. to +30xc2x0 C., preferably below 25xc2x0 C., more preferably in the range of 0 to 10xc2x0 C. The coupling of the diazo compound of an amine of formula (VIII) to a middle component of formula (IX) may be carried out under distinctly acidic to strongly basic pH-conditions, e.g. at a pH in the range of 4 to 12, preferably 5 to 11. The selective coupling of a bisdiazo compound of a diamine of formula (II) or (VI) at one side to the 3-position of a compound of formula (III) is suitably carried out under distinctly acidic conditions, e.g. at a pH in the range of 1 to 5, preferably by adding an aqueous solution of a compound of formula (III) of a nearly neutral pH (e.g. pH in the range of 6 to 8, preferably 6.5 to 7.5) into the acidic solution, in which the diazotization has been carried out and that contains the diazonium compound.\nThe coupling of the monoazodiazo compound to a coupling component Hxe2x80x94B1 or Hxe2x80x94B2 may be carried out at any pH as suitable for coupling to the respective coupling component, preferably at pH values below 10, in particular in the range of 4 to 8.5, more preferably under nearly neutral conditions, in particular at pH-values in the range of 6 to 7.5 and at temperatures, e.g. in the range of xe2x88x925xc2x0 C. to +30xc2x0 C., preferably 0 to 25xc2x0 C. The coupling of a diazo compound of an amine of formula (V) to the 6-position of a compound of formula (IV) is preferably carried out under mildly basic to nearly neutral pH-conditions, e.g. at a pH in the range of 6 to 9, preferably 6.5 to 8 and under temperatures preferably in the range of 0 to 25xc2x0 C. These pH and temperature conditions are also suitable for the selective coupling of the bisdiazocompound of a bisamine of formula (II) or (VI) at one side to the 6-position of a disazo compound of formula (VII). The coupling of the trisazodiazocompound to a coupling component Hxe2x80x94B1 or Hxe2x80x94B2 may be carried under conditions as described above for the coupling of the monoazodiazo compound. By using a suitable base for pH-adjustment (preferably alkali metal hydroxide, e.g. sodium hydroxide or potassium hydroxide, more preferably sodium hydroxide), the dyes may be obtained in a suitable salt form of high hydrosolubility. If desired, they may be obtained in another salt form, e.g. by precipitation by means of acidification and then conversion to another salt form by addition of a suitable base.\nThe synthetized dyes may be purified, if desired, in any manner conventional per se, e.g. by precipitation (e.g. by salting out or by acidification with a strong mineral acid) and filtration and optionally re-precipitation and filtration, and/or by microfiltration and, optionally, ultra-filtration. Where desired, the dye solution may be further purified and/or concentrated by means of suitable membrane filtrations. The purified dye solutions may be dried to dye-powder or -granulate.\nThe dye may, upon isolation or dialysis, be blended with suitable blending agents conventional per se, e.g. alkali metal salts (sodium carbonate, sodium sulphate) non-electrolyte blending agents (mainly oligosaccharides, e.g. dextrine) or/and with surfactants (e.g. as suitable as wetting agents), in particular non-ionic and/or anionic surfactants. There may e.g. be mentioned hydrocarbon sulphonates, sulphonated castor oil, sulphosuccinates, lignin sulphonate, hydrophilic polysaccharide derivatives, ethoxylated alkylphenols or fatty alcohols. If a surfactant is employed the weight ratio of the surfactant to the dye is advantageously in the range of 5:95 to 40:60. If desired, especially if the composition contains a surfactant, as indicated above, it may be formulated with water as a liquid concentrated dye-composition, preferably with a dry-substance content in the range of 10 to 70%, more preferably 20 to 50% by weight, referred to the weight of the composition.\nThe dyes of the invention are, in general, of dark hues, in particular ranging from green to red hues, e.g. bluish black, reddish black or greenish black, and/or deep black.\nThe tetrakisazo dyes of the invention are anionic dyes and are suitable for dyeing substrates dyeable with anionic dyes.\nAny substrate that is dyeable with anionic dyes is suitable as a substrate that may be dyed with the azo dyes of the invention; these include natural and regenerated cellulose, polyurethanes, basically modified high polymers (e.g. basically modified polypropylene), natural or synthetic polyamides or anodized aluminium, in particular, however, leather substrates. The substrate to be dyed may be in any conventional form, e.g. in the form of loose fibers, filaments, yarns, woven or knitted goods, non-woven webs, carpets, half-ready-made and ready-made soft goods and tanned leather or pelts. The dyes may be employed in any desired concentration up to the saturation of the substrate. The dyeing may be carried out by any conventional methods that are suitable for the substrate to be dyed, e.g. by exhaustion or impregnation methods (e.g. padding, spraying, foam application or application with a roller, or printing), preferably from aqueous medium; for synthetic substrates, the dye may optionally also be incorporated into the synthetic mass. Paper may be dyed in the pulp, or after sheet formation.\nThe dyes of the invention are, however, mainly suitable for the dyeing of leather and pelts.\nAny kinds of leather which are conventionally dyed from aqueous medium are suitable, particularly grain leather (e.g. nappa from sheep, goat or cow and box-leather from calf or cow), suede leather (e.g. velours from sheep, goat or calf and hunting leather), split velours (e.g. from cow or calf skin), buckskin and nubuk leather; further also woolen skins and furs (e.g. fur-bearing suede leather). The leather may have been tanned by any conventional tanning method, in particular vegetable, mineral, synthetic or combined tanned (e.g. chrome tanned, zirconyl tanned, aluminium tanned or semi-chrome tanned). If desired, the leather may also be re-tanned; for re-tanning there may be used any tanning agent conventionally employed for re-tanning, e.g. mineral, vegetable or synthetic tanning agents |e.g. chromium, zirconyl or aluminium derivatives, quebracho, chestnut or mimosa extracts, aromatic syntans, polyurethanes, (co) polymers of (meth)acrylic acid compounds or melamine/, dicyanodiamide/and/or urea/formaldehyde resins|. Thus leathers of very high to very low affinity for anionic dyes may be used.\nThe leathers may be of various thicknesses, thus, there may be used very thin leathers, such as book-binder\"\"s leather or glove-leather (nappa), leather of medium thickness, such as shoe upper leather, garment leather and leather for handbags, or also thick leathers, such as shoe-sole leather, furniture leather, leather for suitcases, for belts and for sport articles; hear-bearing leathers and furs may also be used. After tanning (in particular after a re-tanning) and before dyeing, the pH of the leather is advantageously set to values in the range of 4 to 8 (the leather is xe2x80x9cneutralizedxe2x80x9d); depending on the kind of the leather, there may be chosen an optimum pH range, e.g. for grain leather pH values in the range of 4 to 6, for suede leather and split velours and for very thin leathers pH-values in the range of 4.5 to 8, for intermediately dried suede leathers and intermediately dried split velours the pH may range in the scope of 5 to 8. For the adjustment of the pH-value of the leather there may be employed conventional assistants; for tanned leather of acidic character the pH may be adjusted by addition of suitable bases, e.g. ammonia, ammonium bicarbonate or alkali metal salts of weak acids, e.g. sodium formate, sodium acetate, sodium bicarbonate, sodium carbonate or sodium sulphite, of which sodium formate and sodium bicarbonate are preferred. Sodium carbonate and sodium bicarbonate are usable in particular as second bases for the exact adjustment of the superficial pH-value of the leather. Mineral tanned leather may, if desired, also be masked, e.g. with alkali metal formate, oxalate or polyphosphate or e.g. with titanium/potassium oxalate.\nThe dyeing may be carried out in a manner known per se suitably in an aqueous medium and under conventional temperature and pH conditions, in particular in the temperature range of 20 to 80xc2x0 C., preferably 25 to 70xc2x0 C., milder temperature conditions, in particular in the range of 25 to 40xc2x0 C., being preferred for the achievement of deeper penetrations and for the dyeing of wool-bearing skins and furs. The pH-values of the dye-bath may, in general, range broadly, mainly from pH 8 to pH 3; in general the dyeing may be advantageously begun at higher pH-values and concluded at lower pH-values. Preferably the dyeing is carried out at pH-values xe2x89xa74, in particular in the pH-range of 8 to 4, and for the conclusion of the dyeing procedure the pH-value is lowered (e.g. by addition of an acid conventional in the leather dyeing technique such as acetic acid or formic acid) preferably to values in the range between 4 and 3. The dye concentration may range broadly, if desired, up to the saturation degree of the substrate, e.g. up to 5%, referred to the wet weight of the substrate. The dyeing may be carried out in one or more stages, e.g. in two stages, optionally with insertion of charge reversal of the substrate by means of conventional cationic assistants.\nThe dyes of the invention may, if desired, be employed in combination with conventional dyeing assistants, mainly non-ionic or anionic products (in particular surfactants, preferably hydrophilic polysaccharide derivatives, polyoxyethylated alkylphenols or alcohols, lignosulphonates or sulpho group-containing aromatic compounds).\nA fatting may, if desired, be carried out before and/or after the dyeing process, in particular also in the same liquor. For fatting after the dyeing process the fatting agent is advantageously added before the pH of the liquor is lowered, preferably to values between 3 and 4.\nFor the fatting (in particular fat-liquoring) step there may be used any conventional natural animal, vegetable or mineral fat, fat oil or wax, or chemically modified animal or vegetable fat or oil, which include in particular tallow, fish oils, neats-foot oil, olive oil, castor oil, rapeseed oil, cottonseed oil, sesame oil, corn oil and Japanese tallow, and chemically modified products thereof (e.g. hydrolysis, transesterification, oxidation, hydrogenation or sulphonation products), bees-wax, Chinese wax, carnauba wax, montan wax, wool fat, birch oil, mineral oils with boiling range within 300 and 370xc2x0 C. (particularly the so-called xe2x80x9cheavy alkylatesxe2x80x9d), soft paraffin, medium paraffin, vaseline and methyl esters of C14-22-fatty acids; and synthetic leather fatting agents, including esters, in particular partial esters of polybasic acids (e.g. phosphoric acid) with optionally oxyethylated fatty alcohols. Of the above mentioned the methyl ester, the sulphonation products and the phosphoric acid partial esters are particularly preferred. By the term xe2x80x9csulphonationxe2x80x9d for the fatting agents, there is meant generally the introduction of the sulpho group including also the formation of a sulphato group (=xe2x80x9csulphatingxe2x80x9d) and the introduction of a sulpho group by reaction with a sulphite or SO2 (=xe2x80x9csulphitingxe2x80x9d).\nA conventional leather softener, in particular a cationic leather softener may, if desired, be applied in a final step, particularly if fatting has been carried out with a sulphonated fat-liquoring agent.\nThe treated substrate may then be further treated in conventional manner, e.g. rinsed or washed, drained, dried and cured.\nAccording to the invention there may be obtained azodyes that display, even with a relatively low number of hydro-solubilizing substituents in W1, W2, D, B1 or B2, a high solubility in water, especially where the anionic dyes are in alkali metal salt form; they are distinguished by their stability to electrolytes (in particular inorganic ions), specifically also to bases and acids, and are also distinguished, especially on leather, by their build-up and a high degree of insensitivity to variations of the affinity of the leather towards anionic dyes, very level dyeings of outstanding penetration and high colour-yield being obtainable. In these leather-dyeing properties the dyes of formulae (Ic) and especially (Ib) excel over those of formula (Ia) and are, thus, preferred. According to a particular feature of the invention and with a view to these properties, they may advantageously be employed also in a form substantially free of other elektrolytes, as obtainable by dialysis.\nThe dyeings, particularly on leather, have excellent fastness properties, for example wet-fastnesses, fastness to rubbing, light-fastness and stability to PVC-migration. They are readily combinable with other dyes, in particular such with similar tinctorial behavior. There may be obtained very level, intense, fine dyeings, grain side and velours side being very evenly dyed, the shade of the dyeings on different kinds of leather being equal or very similar; in admixture with corresponding compatible dyes with which the dyes of the invention are combinable, in particular with the above trisazo dyes of formula (X), there may also be obtained very intense and regular dyeings of high yield and optimum fastnesses. By the choice of the substituents some of the properties of the dyes (e.g. solubility, shade, build-up, penetration, levelness etc.) may be varied accordingly. The polyazodyes of the invention, especially the preferred ones, are also well combinable with trisazo or other tetrakisazo dyes, there being obtainable very homogeneous shades and penetrations and regular dyeings of optimum fastnesses, also light fastness."}
{"text": "“Digital goods” is a generic label for electronically stored or transmitted content. Examples of digital goods include images, audio clips, video, multimedia, software, and data. Digital goods may also be called a “digital signal,” “content signal,” “digital bitstream,” “media signal,” “digital object,” “object,” and the like.\nDigital goods are often distributed to consumers over private and public networks—such as Intranets and the Internet. In addition, these goods are distributed to consumers via fixed computer readable media, such as a compact disc (CD-ROM), digital versatile disc (DVD), soft magnetic diskette, or hard magnetic disk (e.g., a preloaded hard drive).\nUnfortunately, it is relatively easy for a person to pirate the pristine digital content of a digital good at the expense and harm of the content owners—which includes the content author, publisher, developer, distributor, etc. The content-based industries (e.g., entertainment, music, film, etc.) that produce and distribute content are plagued by lost revenues due to digital piracy.\nModern digital pirates effectively rob content owners of their lawful compensation. Unless technology provides a mechanism to protect the rights of content owners, the creative community and culture will be impoverished.\nWatermarking\nWatermarking is one of the most promising techniques for protecting the content owner's rights of a digital good. Generally, watermarking is a process of altering the digital good such that its perceptual characteristics are preserved. More specifically, a “watermark” is a pattern of bits inserted into a digital good that may be used to identify the content owners and/or the protected rights.\nGenerally, watermarks are designed to be completely invisible or, more precisely, to be imperceptible to humans and statistical analysis tools.\nA watermark embedder (i.e., encoder) is used to embed a watermark into a digital good. A watermark detector is used to extract the watermark from the watermarked digital good. Watermark detection is performed in real-time even on small devices.\nThose of ordinary skill in the art are familiar with conventional techniques and technology associated with watermarks, watermark embedding, and watermark detecting.\nWatermarks have limitations. They may be used to designate a digital good as protected and, perhaps, to indicate that a license is necessary to legally use the digital good. However, since watermarks are identical in all copies of a digital good, a digital pirate can reproduce the original content of a marked copy by breaking the watermark at a single watermark detector, for example by extracting the detection key and using it to find the watermark and remove it or jam it.\nTherefore, others may use the original content without the watermark; thus, without the content owner receiving the appropriate compensation. This is generally called “break once run everywhere” or BORE.\nFurthermore, to individualize a particular copy of a digital good (or a particular system that will use that good) with watermarks, we need to augment it with a technology called “fingerprinting”.\nConventional Fingerprinting\nConventional fingerprinting (i.e., “classic fingerprinting”) refers to techniques that involve uniquely marking each copy of a particular digital good, and associating each uniquely marked copy with a “classic fingerprint.” That classic fingerprint is associated with or assigned to a particular entity (e.g., person, business, media player, or smart card) to which the copy is distributed.\nIf unauthorized copies of the uniquely marked copy are made, the fingerprint can be traced back to the original entity to which the copy was initially distributed. In other words, classic fingerprinting technology may be used to trace piracy to its origin.\nAs an example, consider a printed map. When a mapmaker produces a map, they may want to ensure that those individuals to whom the map is distributed do not make unauthorized copies of the map and distribute them to others. One way that the mapmaker might protect his maps is to introduce a different trivial error (e.g., a non-existent street) into each of the copies of the map that are distributed. Those different trivial errors are fingerprints. Each fingerprint is then associated with an individual to whom the map is distributed. By associating each different fingerprint with a different individual, if and when unauthorized copies of that individual's copy are uncovered, they can be traced back to the original individual by virtue of the unique fingerprint that the map contains.\nUsing embedding methods similar (but not identical) to watermarking, the fingerprint is embedded into a digital good. If we want to achieve both prevention and “after the fact” tracing, a combination of the fingerprint and watermark are embedded into a digital good.\nVery powerful machines that can devote significant resources to the process of detecting a fingerprint typically perform fingerprint detection. If necessary, a fingerprint detector can have access to the original unmarked digital good, using it to improve the likelihood of success in detecting the fingerprints—even from content modified by malicious attacks.\nClassic Fingerprint=Unique Entity Identifier (UEid)\nAlthough the term “fingerprint” is commonly understood by those of ordinary skill in the art, the terms “classic fingerprint” or “unique entity identifier” (UEid) may be used hereinafter to refer to this conventional technology (and its unique marks). This is done to avoid confusion with the use, herein, of “fingerprinting” in the other sections of this document (i.e., sections other than the “Background” section). In those other sections, the term “fingerprinting” may refer to a similar, but distinctly different technology.\nCollusion\nOne problem with fingerprinting can arise when two or more entities collude. Their purpose for doing this may be to discover, modify, or remove their fingerprints and/or the embedded watermark. Those that attempt to collude are called “colluders.” A group of colluders who attempt to collude are part of a “collusion clique.”\nReturning to the map example for illustration, collusion occurs when two or more individuals get together and compare their maps. They can, given enough time, ascertain their unique fingerprints by simply looking for the differences between their maps. If they can ascertain their fingerprint, they can alter it and therefore possibly avoid detection.\nWith the advent of the Internet and electronic distribution, fingerprinting digital goods for purposes of detecting or deterring unauthorized copying has become particularly important. As in the above map example, collusion by different individuals in the digital context can pose challenges to the owners and distributors of digital goods.\nConventional Fingerprinting/Watermarking Systems with Collusion Resistance\nExisting conventional fingerprinting/watermarking systems have some capability for collusion detection. However, the protection offered by these systems is limited.\nFor example, Ergun et al. have proved that no conventional fingerprinting system can have a better asymptotical collusion-resistance than: O((N/log(N))1/2)—where O indicates “order of magnitude” and N is the size of the marked digital good. For example, the best fingerprinting system today, “the Improved Boneh Shaw System” achieves for a typical two hour movie a collusion resistance of only 40 users. This system, just as the original “Boneh Shaw Fingeprinting System” has collusion resistance in the order of O(N1/4).\nThe derivation of the upper bound on fingerprinting mechanisms by Ergun et al. considers embedding distinct fingerprints per copy of a digital good and models collusion attacks as averaging of copies with additive noise. Aspects of their work are described in an article entitled “A Note on the Limits of Collusion-Resistant Watermarks,” authored by Ergun, Kilian, and Kumar, appearing in Proc. Eurocrypt, 1999.\nFor example, another conventional fingerprinting system (the “Boneh-Shaw, or B-S system”) defines a lower bound on collusion-resistant fingerprinting: O(N1/4). Assuming that the marked digital good is a typical music clip, the lower bound of the number of colluders necessary to thwart this conventional system is in the neighborhood of 4. The B-S system is a fingerprinting system that attempts to overcome the problem of collusion when fingerprinting digital goods. Aspects of the B-S system are described in an article entitled “Collusion-Secure Fingerprinting for Digital Data” authored by Boneh and Shaw, appearing in IEEE Transactions on Information Theory, Vol. 44, No. 5, Sep. 1998.\nThose of ordinary skill in the art are familiar with conventional techniques and technology associated with classic fingerprinting, classic fingerprinting embedding, and classic fingerprinting detecting.\nAlthough the conventional fingerprinting systems provide some protection against collusion, that protection is only effective when the number of colluders is relatively small. Consequently, the confidence level that a marked digital good is free from the effects of collusion is not high.\nAccordingly, there is a need for a new watermarking/fingerprinting technology that is more collusion resistant. A new technology is needed that increases the protection that is provided by fingerprinting (and watermarking) to detect colluders even when their numbers are large. If that numbers is several orders of magnitude greater than the conventional, then the confidence level—that a marked digital good is free from the effects of collusion—would be very high indeed.\nMoreover, there needs to be a more effective technique to identify that a digital good has had its mark removed and who removed that mark. That way, piracy crimes can be more effectively investigated."}
{"text": "In some types of partial middle ear implantable (P-MEI) or total middle ear implantable (T-MEI) hearing aid systems, sounds produce mechanical vibrations within the ear which are converted by an electromechanical input transducer into electrical signals. These electrical signals are in turn amplified and applied to an electromechanical output transducer. The electromechanical output transducer causes an ossicular bone to vibrate in response to the applied amplified electrical signals, thereby improving hearing.\nAn electromechanical transducer used for the purpose of vibrating or sensing from any or all elements of the ossicular chain may be mounted in or near the middle ear. The transducer is generally contained in a housing or enclosure, forming a driver or sensor assembly that facilitates the placement of the transducer within the middle ear.\nIn previous designs, applicant has noticed unwanted resonances within the audible frequency range, which can be disconcerting to the person having the implant.\nWhat would be desirable are transducers which more accurately convert the electrical signal received into vibrations which can be coupled to the ossicular chain."}
{"text": "The present disclosure generally relates to information processing and, in particular, storing user data on servers.\nMany service providers (e.g., cloud computing services, data sync services, file hosting services, etc.) enable user data to be stored at one or more servers. For example, an application synchronization service may allow a user to sync one or more applications so that user data associated with one application on one client system may be transferred and used on another application on another client system. Synchronization of the applications on the two systems may be achieved by storing user data on a data storage server.\nIn order to protect user data and the privacy of users, software on a first client system may generate an encryption key based on a passphrase (e.g., the user's password or another user supplied string) and encrypt the user data using the encryption key before it is transmitted to a data storage server to be stored. The user may download the encrypted data from the data storage server on a second client system. However, in order to access the encrypted data on the second client system, the second client system must have access to a key to decrypt the encrypted data. In some cases, the second client system may generate a key using the same user provided passphrase or password and decrypt the encrypted data using the key. However, the user may not be able to generate the key to decrypt the encrypted user data if the user forgets the passphrase (or if the user changes his password)."}
{"text": "The invention relates to an illumination module for illuminating an object to be examined, which module comprises an exit window for radiation generated in the module, a module axis which extends perpendicularly to the central part of the exit window, and a number of radiation-generating elements which are arranged in at least two illumination rings, which illumination rings have a midpoint, which is situated on the module axis, and are individually switchable, the radiations of the different illumination rings passing through the exit window at different angles with respect to the module axis.\nThe invention also relates to an optical device for the observation of an object, which device is provided with such an illumination module, and to an apparatus for placing electronic components on a support, which apparatus is provided with at least one such optical device.\nAt present, said apparatus is often used to automatically, very accurately and very rapidly place electronic components on a support, such as a printed circuit board (PCB). Said apparatus comprises one or more robots, provided with a suction pipette, and an equally large number of placement stations associated with the relevant robots. The suction pipette receives a component to be placed from a component input device, whereafter the associated robot transports the suction pipette and the component to a position on a support which, at that instant, is situated in the placement station. Before the component can be placed on the support, it has to be identified and inspected for irregularities, and it must be checked whether the component is correctly oriented with respect to, for example, a system of coordinates of the robot. To this end, the placement station is provided with an optical observation device for observing the component held by the suction pipette, which observation device is alternatively referred to as vision module. This device comprises a camera by means of which an image of the component is converted into electric signals which are electronically processed. A good observation device should meet the following requirements: the image must be free of distortion, the observation device should also be capable of imaging large components, the aperture angle of the beam with which the component is illuminated must be relatively small, and, for specific applications, the angle of incidence of the illumination beam on the component must be as large as possible. To properly observe the component, it is necessary for this component to be illuminated in the right way. To achieve this, the observation device comprises an illumination module.\nSuch an illumination module is described in U.S. Pat. No. 4,893,223. The known illumination module comprises a large number of light-emitting diodes, also referred to as LEDs, which are mounted on the inner surface of a concave plate and arranged in rings or angle sectors. Since the components to be placed by means of the placement machine may demonstrate different shapes and different surface conditions, the illumination module should preferably be programmable, i.e. it should be possible for groups of LEDs, such as the LEDs of a ring or the LEDs of a comer sector, to be switched on individually, so that the angle at which the object is illuminated and/or the part of the object that is illuminated can be adjusted. The illumination module described in U.S. Pat. No. 4,893,223 is used to illuminate the free end, which is provided with a solder ball, of the connection pins of a component, which is commonly referred to as surface mounted device (SMD). The observation device wherein this illumination module is employed, makes it possible to inspect only the end of one pin of the SMD. To inspect a subsequent pin, the observation device and the SMD are moved with respect to each other. Thus, the illumination module described in U.S. Pat. No. 4,893,223 has a small illumination field, and the angle between the beam of the LEDs and the normal to the surface of the SMD is relatively small. The beams supplied by the LEDs have a relatively small aperture angle, of the order of 10xc2x0. The illumination module in accordance with U.S. Pat. No. 4,893,223 is difficult to manufacture because all LEDs must be accurately aligned to make sure that they all illuminate the same small area.\nIt is an object of the invention to provide a programmable illumination module having a wide range of applications, which can be readily manufactured and has minimal dimensions, and which, in addition, can be adjusted for a large illumination field, and which is capable of supplying a very oblique illumination beam. This illumination module is characterized in that all the radiation-generating elements are provided on a flat plate which extends parallel to the central part of the exit window, and the principal rays of all radiation beams emitted by the radiation-generating elements extend parallel to the module axis.\nThe fact that this module is capable of supplying a very oblique illumination beam means that selected radiation, originating, for example, from the outermost illumination ring, includes a very small angle, of for example at least 6xc2x0, with the illuminated surface of the object.\nThe provision, with great accuracy, of the radiation-generating elements on a flat plate is simpler than the provision, with the same accuracy, of said elements on a concave plate. The flat plate may be embodied so as to be a printed circuit board with electric conductors via which the radiation-generating elements are fed.\nA preferred embodiment of the illumination module is further characterized in that opposite each illumination ring there is provided a radiation guide which, in a plane at right angles to the module axis, is ring-shaped in section, the entrance plane of said radiation guide being situated opposite the associated illumination ring, and the shape of the light guide is exit plane situated opposite the entrance plane being adapted to the desired direction of the radiation to be supplied by the associated illumination ring.\nThe use of these radiation guides enables the optimum illumination direction for a specific application to be achieved. The radiation guide associated with a certain illumination ring also makes sure that the radiation from said ring is not mixed with radiation from other illumination rings. Since the radiation from an illumination ring is confined to the associated radiation guide until it exits via the exit plane of this guide, the beams of the radiation-generating elements can have a relatively large aperture angle, for example in the range from 20xc2x0 to 30xc2x0. In addition, the radiation guide ensures, via total internal reflections, that the radiation beam leaving the radiation guide demonstrates a certain degree of homogeneity. In this regard, the length of the radiation guide and the aperture angle of the LEDs of the associated illumination ring play a part.\nTo further increase the homogeneity, the illumination module is preferably further characterized in that the exit surface of all radiation guides, with the exception of the outermost radiation guide, is provided with a diffuse body the surface of which facing away from the radiation guides has such a shape that the radiation beam exiting from a radiation guide through this surface has a desired direction.\nThe radiation guides may be filled with air. Preferably, however, the illumination module is characterized in that the radiation guides, possibly with the exception of the innermost radiation guide, are filled with a transparent synthetic resin.\nIn accordance with a further characteristic, the synthetic resin is polycarbonate or polymethyl methacrylate. These synthetic resins, which are also known by the abbreviations PC and PMMA, respectively, are highly transparent and of good optical quality.\nPreferably, the illumination module is characterized in that scattering particles are integrated in the synthetic resins.\nIf so, these synthetic resins have a diffuse effect, and the relevant radiation guides can contribute to rendering the radiation leaving the illumination module uniform. The different radiation guides may exhibit different degrees of diffusion.\nUnder certain conditions it may be attractive to also fill the innermost radiation guide with a synthetic resin.\nIn accordance with a further characteristic of the illumination module, the outermost radiation guide has a reflective surface for reflecting radiation towards the exit plane of this radiation guide, which reflective surface consists of two parts which enclose an obtuse angle with each other.\nSaid reflective surface divides the radiation beam, which propagates through the radiation guide and which has a certain intensity distribution, into two sub-beams which are positioned at the location of the object such that the original intensity distribution is lost and a more uniform distribution is obtained. In this manner, also the radiation beam supplied by the outermost radiation guide obtains a sufficiently uniform distribution.\nAn embodiment of the illumination module which is attractive for the observation device of a component placement machine is characterized in that it comprises three illumination rings.\nThis illumination module is still compact enough to be fitted in the placement machine, and it also enables good illumination, with sufficient intensity and at different angles, of the components to be achieved.\nFor the radiation-generating elements use can be made of different elements. For example, use can be made of diode lasers. Diode lasers emit a relatively narrow beam of high intensity and the radiation thereof can be readily oriented in the desired direction. For the radiation-generating elements use can also be made of luminescing diodes, provided the elements do not have to be switched rapidly. A preferred embodiment of the illumination module is however characterized in that the radiation-generating elements are formed by LEDs.\nThese LEDs are inexpensive and available in different versions as regards radiation intensity, dimension and aperture angle.\nDependent upon the application of the illumination module, one or more illumination rings may comprise diode lasers and the other illumination rings may comprise LEDs. It is alternatively possible that one and the same illumination ring comprises a number of diode lasers as well as a number of LEDs.\nThe illumination module is preferably further characterized in that at least one of the illumination rings comprises at least two types of LEDs, which emit radiation of different wavelengths, and LEDs of the first type alternating with LEDs of the second type in said illumination ring.\nFor example, all illumination rings may comprise alternately red, green and blue LEDs which can be switched on in different combinations. As a result, the color of the illumination can be optimally adapted to the color of the object to be examined and/or the background thereof. The contrast of the observation device can thus be improved considerably, also if the camera of the observation device is a monochrome camera. It is alternatively possible to use LEDs of only two different colors, or provide LEDs of different colors in only one illumination ring, or provide only one type of LED in one and the same illumination ring; however, the LEDs of the different illumination rings are of a different kind.\nAn illumination module which is particularly suited for use in a component placement machine may further be characterized in that, within the central illumination ring, a transparent plate is provided which encloses an acute angle with the module axis.\nBy means of this plate, it can be precluded that a component which has fallen from the suction pipette situated above the illumination module can fall on the camera situated below the module.\nThis illumination module may be further characterized in that an optical filter is provided on the plate.\nAs a result, this plate may fulfill a second function. The filter may be a color filter by means of which it is precluded that ambient radiation, i.e. radiation not originating from the illumination module, can reach the camera. The filter may also be a grey filter by means of which the intensity of the illumination radiation originating from the component can be adapted to the sensitivity of the camera. This enables the illumination module to be used in combination with different cameras, or to further increase the intensity of the illumination radiation to a level above that of the ambient light so as to reduce the influence of the ambient light on the observation.\nAnother embodiment of the illumination module is characterized in that, within the central illumination ring, a semi-transparent plate is provided at an angle of, in principle, 45xc2x0 with respect to the module axis, and in that an additional radiation source is provided which is arranged such that radiation from this source is incident on the semi-transparent plate and reflected thereby towards the exit window of the illumination module.\nAs the rays of the additional radiation source enclose angles in the range from 0xc2x0 to, for example 20xc2x0 with the module axis, the module can now provide a bright-field illumination for the object in addition to the dark-field illumination which is supplied by the radiation-generating elements of the illumination rings. In the case of a bright-field, or normal, illumination, radiation which is specularly reflected by the object is captured by the imaging optics and the camera, and the object is imaged as a bright object. Details of the object which scatter the illumination radiation, or reflect it outside the imaging optics, can now be observed against a bright background. In the case of a dark field, or oblique, illumination, the object reflects the illumination radiation outside the imaging optics. The scattering details of the object reflect the radiation in the imaging optics and the camera and are brightly imaged against a dark background.\nThe last-mentioned embodiment is preferably further characterized in that the additional radiation source is formed by a number of additional radiation-generating elements arranged along at least one row, in combination with a radiation-conducting plate the entrance face of which is directed towards the additional radiation-generating elements, and the exit face of which is directed towards the semi-transparent plate.\nPreferably, the additional radiation-generating elements are arranged in at least one circle segment at one side of the semi-transparent plate.\nIn this case, the bright radiation field is shaped so as to be in the form of a ring segment the length of which is determined by the position of the midpoint of the semi-circle with respect to the center of the transparent plate.\nIn order to obtain a more two-dimensional foreground illumination, the illumination module is preferably further characterized in that a second radiation-guiding plate is arranged against the exit face of the first radiation-guiding plate at an angle of, in principle, 90xc2x0 with respect to the first radiation-guiding plate, and in that a reflecting roof prism is situated at the location of the transition from the first radiation-guiding plate to the second radiation-guiding plate, so as to reflect, in two opposite directions in said plate, two halves of radiation exiting from the first radiation-guiding plate.\nAs said radiation halves are subject to total internal reflection in the second radiation guide, the beam incident on the semi-transparent plate is elongated in a direction transverse to the direction of the row of additional radiation-generating elements, and this beam additionally demonstrates a more uniform intensity distribution. The roof prism is embodied in such a manner, i.e. the angles thereof are chosen to be such that a part of the radiation leaving the first radiation guide is passed by the prism, so that radiation also exits through that, central, part of the exit plane of the second radiation-guiding plate, which is situated at the level of the roof prism.\nIn accordance with a further characteristic of this illumination module, the roof prism is provided in the second radiation-guiding plate.\nPreferably, however, the illumination module is alternatively characterized in that the roof prism forms part of the first radiation-guiding plate.\nIf so, the prism can be more readily provided and, in addition, the optical properties of the second radiation-guiding plate are not adversely affected.\nPreferably, this illumination module is further characterized in that the second radiation-guiding plate takes the form of a cylinder segment whose cylinder axis extends parallel to the module axis.\nThe outer ray of the second radiation-guiding plate can then be chosen to be such that the upper part of this plate can be slid into the innermost illumination ring, so that, also if the additional radiation-generating elements for the bright-field illumination are employed, the illumination module can remain compact.\nThe additional, radiation-generating elements can alternatively be formed by at least two kinds of elements which emit radiation of different wavelengths, whereby, in a row, elements of the first kind alternate with elements of the other kind(s). In this case, also the color of the bright-field illumination can be adjusted, which may yield a higher contrast.\nThe invention also relates to an observation device for observing an object, which observation device is provided with an illumination system for illuminating the object, and with a camera for receiving an image of the object. The observation device in accordance with the invention is characterized in that the illumination system is an illumination module as described hereinabove.\nFinally, the invention also relates to a component-placement machine for placing electronic components on a support, which machine is provided with:\na frame;\nat least one robot;\na transport system for transporting supports;\na placement head for each robot for placing components on a support, which placement head is secured to an arm of the robot, and\ncomponent-positioning means for each robot for positioning a component which is held by the placement head associated with the robot.\nThis placement machine is characterized in that the observation device is a device as described hereinabove."}
{"text": "In the area of electrical power, a bushing is an insulated device that allows an electrical conductor to pass safely through an (usually) earthed conducting barrier such as the wall of a transformer, circuit breaker or reactor. A condenser bushing is a bushing in which metallic or nonmetallic conducting layers are arranged within the insulating material for the purpose of controlling the distribution of the electric field of the bushing, both axially and radially by capacitive grading. In many cases, bushings may be used with very high voltages such that a failure in the bushing may lead to arcing and a catastrophic failure of the attached device. Accordingly, testing and monitoring of the bushing is very advantageous.\nThe current solutions for on-line transformer monitoring are passive systems that can only provide bushing capacitance and power factor values for the inner condenser section of energized bushings. The capacitance and power factor values for both the transformer and for the outer condenser section of an energized bushing cannot be measured. The embodiments described herein are a significant improvement in the art that address, or at least ameliorate, some of the shortcomings of the current art.\nFIG. 1 illustrates a current design of a bushing monitoring system comprising a condenser-type bushing 10, a passive sensor 20, and a monitor processor 30. High voltage condenser-type bushings 10 may be designed in a wide variety of constructions, and the description herein is not intended to describe all of the possible types of designs. However, the majority of high-voltage condenser-type bushings 10 currently manufactured are constructed with an inner capacitive condenser section 12 referred to as C1, and, an outer capacitive condenser section 14, referred to as C2. These bushing sections C1 and C2 are wound radially around the bushing's central conductor 11. The inner C1 section 12 comprises the bulk of the total condenser 13. The outer C2 section 14 of the condenser 13 normally contains only 5% to 10% of the total number of layers in the complete condenser 13. This type of bushing 10 is normally equipped with an external tap connection 16, which is coupled to a foil layer between the inner C1 section 12 and the outer C2 section 14. The tap 16 may be referred to by several different names such as a test tap, capacitance tap, voltage tap, power factor tap or other designations. For the purposes of this document, we will refer to the tap 16 as a test tap.\nThe inner C1 section 12 of the bushing condenser 13 provides a capacitive structure, which is intended to provide a consistent voltage stress gradient from the energized inner conductor 11 to the outer layer of the inner C1 section 12. The outer layer of the inner C1 section 12 is normally connected to an earth ground 18 when the bushing 10 is in-service. Grounding of the inner C1 section 12 is normally accomplished by attaching a cover onto the test tap 16. The cover electrically connects the test tap 16 to the grounded metal flange of the bushing.\nIn many condenser designs, the outer layer of the outer C2 section 14 is also grounded. In such embodiments, the outer layer of the outer C2 section 14 of the bushing 10 may be comprised of a conductive foil which is connected to the grounded metal flange of the bushing 10. When this type of bushing 10 is in-service and the test tap cover is installed, the C2 condenser section 14 is either directly or effectively earth grounded on the inner and outer layers. However, in some bushings, normally in bushings designed for lower voltages, the outer layer may be ungrounded and separated from the grounded metal flange of the bushing 10 by an insulating fluid or dielectric.\nTypically, the C1 and C2 bushing condenser sections 12 and 14 have capacitance values in the pico-farad range which result in high impedances and low leakage currents through the condenser 13. Typically, the C1 and C2 condenser sections 12 and 14 also have low power factor values that are normally in the 0.2% to 0.5% range as a result of the low power losses in the condenser 13.\nDuring the normal manufacturing process, the C1 and C2 bushing sections 12 and 14 are tested for capacitance and power factor. This ensures the bushings sections 12 and 14 are manufactured without defects. In addition, field tests can be performed after the bushings are installed to detect deterioration or damage to the condenser 10. If the condenser 10 deteriorates sufficiently, the bushing 13 can fail catastrophically, which can cause consequential equipment damage or unsafe conditions for workers in the area. Often, off-line tests and on-line monitoring are performed with the goal to detect such deterioration before a failure occurs.\nBushings are commonly used with transformers or reactors. The transformer or reactor windings have self and mutual capacitances that represent the quality of the insulation for each winding to the grounded parts of the transformer, or between the different windings inside the transformer. Each winding has a capacitance to ground that is a function of the physical dimensions and clearances of the windings as well as the condition of the insulation. For the high voltage (“HV”) winding, this has traditionally been referred to as CH. For the low voltage (“LV”) winding, this has traditionally been referred to as CL. And the capacitance between the HV and LV windings has traditionally been referred to as CHL. Similar designations can be made for three-winding transformers with HV, LV and tertiary voltage (TV) windings to provide CH, CL, CT, CHL, CHT and CLT values.\nOff-line tests, when the equipment is not in-service, have been performed on transformers, reactors and bushings for several years. In order to test a condenser-type bushing's inner C1 section 12, the bushing's test tap cover is removed and test voltages are applied to the center conductor 11 and the resulting current magnitude and relative phase angle are measured from the test tap 16. In order to test the outer C2 section 14, the test tap 16 is energized and the resulting current from the grounded metal flange of the bushing is measured for amplitude and relative phase angle. Off-line tests on higher voltage bushings are normally performed at significantly reduced voltage levels since it is not practical to produce the rated bushing voltages with portable field test equipment. Off-line transformer tests are performed at the same reduced test voltage levels with the voltage applied to the top terminals of the bushings and the resulting current magnitude and phase angle shift are measured to calculate capacitance and power factor values for each winding.\nOn-line bushing monitoring, when the equipment is energized and in-service, has been performed for a shorter time. The current art connects a passive sensor 20 to the test tap 16 and measures the resulting test tap current magnitude and phase angle when the bushing 10 is energized by the power system. As may be seen in FIG. 1, the typical passive sensor 20 included a shunt resistor or capacitor 22 (labelled “A”) and a voltage arrestor 24 (labelled “B”). This method only measures the power frequency current 26 (labelled “IP”) through the inner C1 section 12 of the condenser 13 so only the C1 capacitance and power factor can be determined. Neither the outer C2 section 14 of the condenser 13 nor the transformer winding capacitances can be monitored with the current method.\nAdditionally, transformers with wye winding connections may have neutral bushings that are connected to earth ground either directly or through an impedance, to stabilize the three-phase voltages and limit the phase-ground voltages to safe levels. These neutral bushings can be tested with off-line testing equipment, but neither the inner C1 section 12 nor the outer C2 section 14 of the bushing 10 can be monitored on-line with the current art since the center conductor 11 is effectively grounded and no voltage nor current are produced at the test tap 16."}
{"text": "1. Field of the Invention\nThe present invention relates to a plasma welding method for a golf club head. Particularly, the present invention relates to the plasma welding method of tilting a plasma nozzle a predetermined angle with respect to the golf club head for avoiding obstruction of a neck portion or a hosel of the golf club head. More particularly, the present invention relates to the plasma welding method for welding on a connection line, within a predetermined section, between a striking plate and a club head body.\n2. Description of the Related Art\nReferring now to FIGS. 1 and 2, a conventional golf club head includes a club head body 91 and a striking plate 92 mounted to the club head body 91, as disclosed in TWN Patent Pub. No. 585,792, entitled “golf club head and manufacturing method therefor (II),” TWN Patent Pub. No. I225,421, entitled “golf club head structure,” TWN Patent Pub. No. I226,251, entitled “connecting structure for a striking plate with a golf club head,” and U.S. Pat. No. 6,099,414, entitled “golf club head and method for producing the same” etc. A connection line 93 exists between the club head body 91 and the striking plate 92 which are welded by a suitable welding method. In welding operation, a variety of welding methods, such as tungsten inert gas (TIG) welding, laser welding, plasma welding or other suitable welding methods, may be selectively used.\nThe above plasma welding may be selected to weld the club head body 91 and the striking plate 92, which carries out advantages of high-energy, high-accuracy, high-speed and high-quality welding, and further minimizes or attenuates a heat affected zone of material in a welding process.\nIn plasma welding operation, a movable plasma nozzle 94 is disposed on a slide track (not shown) and aligned with the connection line 93 existing between the club head body 91 and the striking plate 92. Subsequently, the plasma nozzle 94 is moved downward to an operating level above the connection line 93, and welded along the connection line 93 by an automatic welding process. The club head body 91 has a neck portion 91 and a hosel 912 connected thereto. It should be noted that the neck portion 911, however, is located between the hosel 912 and the club head body 91 and extended upward from a level of the striking plate 92. The neck portion 911 may unavoidably obstruct or interfere in a runway of the plasma nozzle 94 in the plasma welding since the neck portion 911, to a certain extent, is too close to the connection line 93. Consequently, the plasma nozzle 94 running at the operating level cannot pass through above a section of the connection line 93 adjacent to the neck portion 911 of the club head body 91 due to interference with the neck portion 911.\nAs explained above, the automatic plasma welding process is only suitable for welding the other section of the connection line 93 away from the neck portion 911 of the club head body 91. Meanwhile, other suitable welding methods, a TIG welding method for example, may be used to manually weld the section of the connection line 93 adjacent to the neck portion 911 of the club head body 91. Such welding method may cause poor quality, increased deformation and low efficiency of welding on the connection line 93 adjacent to the neck portion 911 of the club head body 91. Accordingly, such practice, however, may limit the welding process to carry out to be a complete automatic welding process. Hence, there is a need for improving the plasma welding process for the golf club head.\nThe present invention intends to provide the plasma welding method of tilting a plasma nozzle a predetermined angle with respect to a golf club head for avoiding obstruction of a neck portion or a hosel of the golf club head in such a way to mitigate and overcome the above problem. This permits the plasma nozzle passing through a connection line adjacent to the neck portion of the golf club head in the plasma welding process. Accordingly, the plasma welding method carries out a complete automatic plasma welding process so as to improve the welding efficiency and quality."}
{"text": "1. Field of the Invention\nThe present invention relates to internal combustion engines and, more particularly, to an engine power enhancement/pollution control system utilizing vortex separation to enrich the air input with oxygen.\n2. Description of the Related Art\nThe use of oxygen enriched air to reduce the content of toxic compounds from exhaust gases of automotive engines is known in the prior art. Nakajima et al., U.S. Pat. No. 3,817,232, for example, discloses the delivery of denitrified air, containing oxygen in a major proportion to a carburetor of an internal combustion engine. The disclosed apparatus, however, is applicable to an internal combustion engine only after major modifications in the engine structure. Moreover, the patented structure requires the use of two pumps, forming an integral part of the air intake system for the engine, along with an air denitrifying unit. The latter operates by using a nitrogen impermeable membrane, for example, or a specified molecular sieve formed of pulverized zeolite.\nSuch a structure is complex, expensive, requires major engine modification, and is thus now easily adaptable for use with older cars, subsequent to production and sale. McKerahan, U.S. Pat. No. 1,339,211, discloses the use of a rotary concentrator for delivering oxygenated air at its output, in order to obtain a fuel saving by more complete combustion in smelting furnaces, blacksmith fires, steam boilers, gas engines and the like.\nSimilarly, U.S. Pat. No. 4,351,302, issued to D. H. Brettler, discloses an oxygen concentrator which utilizes a frustoconical structure for reducing the nitrogen content and enriching the oxygen content of air provided for use in an internal combustion engine. The Brettler device utilizes density differences, i.e. separation of heavier from lighter components of the incoming air, to separate the oxygen and nitrogen components thereof. Essentially, the Brettler device serves as a centrifuge. Its separation efficiency is quite low.\nSeparation processes are one of the most widely used industrial processes, especially in chemical and petrochemical industries. Such separation processes are also one of the most expensive industrial processes. They require costly capital investments such as distillation columns and high utility expenditures for both heating and cooling. Currently, the fractional distillation process is almost exclusively used for separation of species with high process rates. Other low flow rate processes, such as membrane, ion exchange, etc. are available with limited applications.\nSeparation of hot gases and cold gases by vortex separation methods is known. A single phase (gas) and single component (pure nitrogen, for example) vortex tube is known as the Ranqe-Hilsch tube illustrated in FIGS. 1 and 1b (Prior Art), designated generally as 1. In the tube, compressed gas is introduced through the nozzle 2 which is directed tangentially to the tube. The gas generates a vortex as it travels into the center of the tube and propagates through the tube.\nAs the gas travels into the center of the vortex, the velocity of the gas increases due to the pressure gradient. The velocity is eventually reduced as the gas travels further into the tube center because the viscosity of the gas slows the fluid.\nWhen the gas slows down in the center of the vortex, the gas has to surrender its kinetic energy. Kinetic energy in the vortex tube is released by transferring the energy from the inner to outer vortex. Thus, the energy separation occurs between the inner and outer vortex causing the temperature differential between the cold inner and hot outer vortex. The result of such an energy separation is that the hot gas 3 is directed in a first direction and cold gas 4 is directed in an opposite second direction.\nAlthough the use of this Ranqe-Hilsch tube 1 is useful for refrigeration applications, it is not useful for providing separation of chemical species."}
{"text": "The invention relates to a vaporization burner for a heating device operated by means of a liquid fuel, particularly for a motor vehicle heating device.\nGerman Patent No. 32 33 321 and corresponding U.S. Pat. No. 4,530,658 disclose a vaporization burner, of the kind mentioned above, which has a combustion chamber within a combustion pipe. The combustion chamber contains an absorbent body disposed on a carrier to which a liquid fuel is supplied. At least part of the side of the carrier facing the combustion chamber is exposed. Furthermore, a combustion air inlet arrangement is provided to prepare a combustion mixture of fuel and air that is ignited by means of an ignition device and burned in the combustion chamber.\nGerman Offenlegungsschrift No. 34 23 940 describes another variation of a vaporization burner for a heating device of the kind initially mentioned, in which the ignition device is an elongated glow plug which transverses the carrier and the absorbent body, projecting above it and extending into the combustion chamber. Still another vaporization burner is disclosed in German Patent No. 32 33 319 and corresponding U.S. Pat. No. 4,538,985, in which a glow ignition device is formed by an essentially flat glow heating body, made of heatable filaments, disposed parallel to the vaporizing surface, immediately in front of the vaporizing body or in the vaporizing body.\nIn all variations of the vaporization burners described above, the glow device extends into the combustion chamber, thereby exposing it to the high temperatures prevailing therein, causing a reduction in the service life of such a device. Moreover, the fact that the glow device extends into the combustion chamber makes it an impedance to the flow conditions within the combustion chamber itself.\nIt is, therefore, a primary objective of the invention to create a vaporization burner for a heating device operated with liquid fuel which overcomes the disadvantages described above, and which provides an extended service life for the ignition device, as well as an improvement with respect to the flow conditions in the combustion chamber.\nIn accordance with the preferred embodiments of the invention, a vaporization burner for a heating device operated with liquid fuel, especially a motor vehicle heating device having a combustion pipe, a combustion chamber in which an absorbent body is disposed on a carrier, to which liquid fuel is supplied and at least part of the side of the carrier facing the combustion chamber is exposed, a combustion air supply device, and an ignition device for the combustion mixture, is provided wherein the ignition device is arranged behind the carrier, i.e., at the side thereof facing away from the combustion chamber. Furthermore, the carrier has at least one opening extending through it that is covered by the absorbent body, in addition to at least one through-opening in the exposed part of the carrier.\nAccordingly, in the vaporization burner of the invention, the ignition device is disposed outside of the actual combustion chamber, i.e., behind the carrier for the absorbent body which receives liquid fuel. As a result, the ignition device does not impede the flow conditions in the combustion chamber, thereby resulting in a more favorable combustion process. Moreover, such an arrangement prevents the ignition device from being directly exposed to the high temperatures in the combustion chamber and/or directly to the flame, resulting in a considerably longer service life for the ignition device than has been previously possible. In other words, a vaporization burner is obtained which operates with the utmost reliability, even over a long period of time.\nThe presence of at least one opening in the part of the carrier that is covered by the absorbent body permits fuel vaporized by the absorbent body to pass directly into the area of the ignition device so that, in combination with the combustion air supplied, an ignitable mixture can be prepared directly in the area surrounding the ignition device. When this mixture is ignited by means of the ignition device, a pilot flame is generated which extends through the through-opening in the exposed part of the carrier. This pilot flame then ignites the fuel-air mixture on the front side of the carrier and the absorbent body (i.e., at the side facing the combustion chamber), thereby initiating the combustion process in the combustion chamber. As is customary, the ignition device is turned off once the combustion is proceeding on its own in the combustion chamber.\nPreferably, the ignition device is arranged in a chamber defined by an ignition housing at the rear of the carrier and such housing has at least one combustion air supply inlet which terminates in the ignition chamber formed thereby. It is also preferable that this combustion air supply inlet be opposite the opening in the carrier that is covered by the absorbent body so that the combustion air supplied to the ignition chamber is locally concentrated around the ignition device to promote mixing, in the ignition chamber, of the air with the fuel vaporized by the absorbent body, for the purpose of quickly obtaining a mixing ratio favorable to the ignition process.\nPreferably, one housing wall of the ignition chamber is formed by the rear side of the carrier, resulting in a minimum space arrangement for the ignition device with the associated chamber. Further, such a construction facilitates placement of the ignition device in the closest possible proximity to the carrier supporting the absorbent body so that it can be rapidly heated when the ignition device is turned on, thereby promoting fuel vaporization.\nIn accordance with another aspect of the invention, the ignition device is surrounded by the housing from its free end up to its mounting thread. In this manner, cooling of the ignition device by the combustion air and passing of air into the chamber around the ignition device is prevented.\nFor reasons of functional efficiency, it has proven beneficial to divide the housing into an outer housing part, representing a separate component and containing the screw thread portion of the ignition device, and an inner housing part which is located in the vicinity of the carrier. Such a design aids in overcoming difficulties caused by variations in manufacturing tolerances, and, more particularly, facilitates direct contact of the outer housing with the ignition device, whereby the ignition device is threadingly engaged with this outer housing part. The inner housing part, on the other hand, has an enlarged interior space to accommodate the premixing of the fuel vapor and combustion air for generating the pilot flame. In order to ensure that the combustion air enters the ignition chamber only via the one combustion air inlet, and to further ensure that no so-called \"false air\" reaches the chamber, the upper and lower housing parts are sealed at their junction.\nFor the purpose of reducing sealing difficulties, the vaporization burner and the ignition device, in accordance with another advantageous variation of the invention, are enclosed by a cup-shaped housing part which is connected with the combustion pipe in the chamber by a single flange connection. Also, to avoid sealing problems, the flange is disposed in the plane of a ring of air openings for the supply of combustion air to the combustion pipe.\nAppropriately, the flange of the combustion pipe defines a slot-shaped opening for the combustion air, which air is directed to the vaporization burner at a distance from the absorbent body on the carrier, in a direction into the combustion chamber.\nIn an advantageous manner, the ignition device forms an integrated unit with the cup-shaped housing part, which unit, appropriately, also includes the carrier with the vaporization element so that, if required, all of the elements essential to the functioning of the vaporization burner can be quickly and simply exchanged as a combined unit by removal and replacement of the cup-shaped housing part.\nIn a further aspect of the invention, the ignition device of the vaporization burner is arranged approximately parallel to the carrier surface. To this end, the ignition device, preferably, is formed by an elongated glow plug.\nIn accordance with an alternative variation of the invention, the housing surrounding the ignition device can form an integrated unit with the ignition device, and this unit can be removable in a radial direction without disassembly of the heating device. Such an arrangement has an advantage with respect to ease of maintenance.\nThese and further objects, features and advantages of the present invention will become more obvious from the following description when taken in connection with the accompanying drawings which show, for purposes of illustration only, several embodiments in accordance with the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a fixing device that causes a toner image on a paper to be fixed to that paper at a nip area formed between a fixing member and a pressing member, and to an image forming apparatus comprising the same.\n2. Description of the Related Art\nA fixing device of -an image forming apparatus generally includes a fixing roller, or a fixing belt (fixing member) wound onto the fixing roller, and a pressing roller (pressing member) that contacts the fixing member to form a nip area between itself and the fixing member. A toner image on a paper that is conveyed in is thermally fixed onto the paper in the nip area. With the fixing device, there are cases where a paper that has undergone the fixing process is discharged from the nip area while being stuck to the fixing member, and in these cases problems are caused such as paper jam.\nTo eliminate such problems, in a first technology for example, the fixing device is provided with a separation mechanism, which is configured so as to be capable of stripping from the fixing member the paper that has stuck to the fixing member after the fixing process. The separation mechanism has a single separation claw (separation member) that extends along the longitudinal direction of the fixing roller and that is arranged facing the fixing roller. A gap between the separation claw and the fixing roller is set to a predetermined magnitude and due to this the separation claw touches the paper that has stuck to the fixing roller and is able to strip the paper from the fixing roller.\nIn the first technology, a spacer that is installed near a leading edge of the separation claw is used a means for setting the gap to a predetermined magnitude. The spacer is set in position so as to abut the core of the fixing roller, and therefore a gap of the predetermined magnitude is formed between the leading edge of the separation claw and the peripheral surface of the fixing roller.\nHowever, only a single separation claw is used in the separation mechanism of the first technology, and therefore due to unevenness in the dimensional precision of the separation claw itself, it is difficult to set the gap to the predetermined magnitude. As a result, it is difficult to strip away the paper that has stuck to the fixing roller and there is a risk that paper jam will occur."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for detecting abnormal execution of program, which detects abnormality in execution state of a program, and particularly to an apparatus for detecting abnormality in execution state of a control program in an electronic control unit for automobile.\n2. Description of the Related Art\nFIG. 3 is a block diagram showing the conventional structure of a door lock control system to which the present invention is applied. This control system is constructed so that a data processor (an electronic control unit 3 for automobile) controls an electric instrument (a door lock motor 5) to be controlled according to input from input means (a door lock switch 1), and a CPU 3a of the electronic control unit 3 controls the door lock motor 5 based on a predetermined control program.\nIn the control program, generally, a series of predetermined operations (predetermined process) is repeated to perform a control operation, as shown in FIG. 4. In the series of the operations, a plurality of sub-routines (sub-routines SR1 to SR4) are executed. Here, in the sub-routine SR1, reading of input from the door lock switch 1 is performed; in the sub-routine SR2, filtering is performed in which whether the input from the door lock switch 1 is not a noise but right is detected; in the sub-routine SR3, control determining operation for detetermining control of the door lock motor 5 in accordance with the input from the door lock switch 1 is performed; and in the sub-routine SR4, the detetermined control is output to the door lock motor 5 and an outputting operation for drive-controlling the door lock motor 5 is performed.\nIn this electronic control unit 3, as shown in FIG. 4, when the electronic control unit 3 is actuated and an initial operation is performed in step ST1, sequential steps ST2 to ST5 are repeated, whereby the sub-routines SR1 to SR4 are repeatedly executed. As a result, the control operation of the door lock motor 5 is performed based on the input from the door lock switch 1.\nIn the CPU 3a of the electronic control unit 3 for executing this control program, during execution of the control program, the memory of the program counter which indicates an executing position of the program (current position) is often disturbed by abnormality of a ROM or effects of noise, so that each of the sub-routines SR1 to SR4 is not normally completed and the executing position of the program can be moved to the unexpected wrong position. In this case, there is possibility that runaway of the CPU 3a (electronic control unit 3) is caused.\nFor example, as shown in FIG. 5, in a series of the operations, in case that step ST3 which must execute the sub-routine SR2 is not executed but step ST4 is executed after the sub-routine SR1 is completed in step ST2, as the sub-routine SR2 is not executed, difference occurs among the various data in the control program and a wrong operation or runaway can occur. Further, to the contrary, there is also a case that each of the sub-routines SR1 to SR4 is wrongly repeated and unnecessarily executed in a series of the operations, which may also cause the wrong operation or the like.\nHowever, in the above-described conventional art, in case that any one of the sub-routines SR1 to SR4 to be executed is not executed or in case that any one of the sub-routines SR1 to SR4 is unnecessarily executed repeatedly in a series of the operations to be repeatedly executed of the control program, measures to meet the respective cases are not taken. Therefore, there is a problem that the wrong operation and runaway of the electronic control unit 3 caused by not suitably executing the predetermined sub-routines SR1 to SR4 to be executed thoroughly cannot be prevented.\nIn view of the above problem, an object of the present invention is to provide an apparatus for detecting abnormal execution of program, which can previously prevent the wrong operation and runaway of a data processor caused by not suitably executing sub-routines to be executed thoroughly in a predetermined process of the control program.\nIn order to achieve the above object, according to the invention, there is provided an apparatus for detecting abnormal execution of program, which detects abnormality in execution state of program in a data processor that performs a predetermined data processing by executing a program having a plurality of sub-routines to be executed in a predetermined process, comprising: memory means which is provided in the data processor and stores a set value representing the number of the sub-routines to be executed in the process of the program; counting means which is provided in the data processor and counts the number of the sub-routines that have been actually executed in the process of the program; and detecting means which is provided in the data processor, compares the set value stored in the memory means with a count value of the counting means in the last step of the process and then detects that the execution state of the program is abnormal in case that the two values are different from each other."}
{"text": "1. Field of the Invention\nThe invention pertains to the field of indexable cutting inserts. More particularly, the invention pertains to an indexable, double sided, nonagon cutting insert.\n2. Description of Related Art\nModern high-performance cutting tools use replaceable and typically indexable inserts owing to the high cutting speeds and feeds supported by the superior insert materials. Common materials for inserts include tungsten carbide, polycrystalline diamond and cubic boron nitride.\nIndexable inserts use a symmetrical polygonal shape, such that when the first cutting edge is blunt they can be rotated or flipped over, presenting a fresh cutting edge which is accurately located at the same geometrical position. Geometrical repeatability saves time in manufacturing by allowing periodical cutting edge renewal without the need for tool grinding, setup changes, or entering of new values into a CNC program.\nCommon shapes of indexable inserts include square, triangular and rhombus (diamond) providing four, three and two cutting edges, respectively, on each side of the insert. A double-sided or invertible square insert, for example, can be flipped over to provide eight cutting edges.\nThe number of cutting edges is directly related to the cost per edge of the cutting insert. The more cutting edges that are available, the more the cutting insert has value. Thus, it is desirable to provide a cutting insert with multiple cutting edges."}
{"text": "1. Field of the Invention\nThe present invention relates to a printing system and more specifically to a configuration of obtaining the status of a printing apparatus.\n2. Description of the Related Art\nConventionally, there is known a printing system with a printer driver which controls printing and a printer utility which controls printer maintenance, both of which are installed in a host computer. In this type of printing system, a user generates printing data by a printer driver and sends the printing data to a printer. On the other hand, there is known a printer utility which performs maintenance such as cleaning of the head of an inkjet printer or adjustment of the head's position, as well as checking the remaining amount of ink or battery in the printer (Japanese Patent Laid-Open No. 10-301728). The printer utility is executed by the host computer.\nThere are cases in which a printer utility and printer driver are or are not installed in a coupled manner in a printer system. For example, in certain operating systems such as operating systems of a Macintosh (Registered Trademark), the printer utility and printer driver are installed in the host computer as different modules. In such a case, when a user starts a printer utility and selects a desired function, the printer utility directly performs sending and receiving of data, and executes each type of process.\nIn the case that the printer utility and printer driver are installed as separate modules in the host computer, the printer driver opens the printer interface upon starting print processing, and performs sending and receiving of data with the printer. When print processing is finished, the interface is closed. During the time between the opening and closing of the interface, the interface is in an occupied state by the printer driver, and other processes cannot access the printer. Therefore, when the printer driver is executing print processing, the printer utility cannot communicate with a printer, and, for example, cannot obtain the status of the printer. For this reason, the user cannot check the remaining amount of ink or battery of the printer while in the middle of printing, nor the printer's operation mode."}
{"text": "The claimed subject matter relates generally to computer communication and, more specifically, to techniques that provide legacy applications with access to a cluster infrastructure."}
{"text": "1. Field of the Invention\nThe present invention relates in general to video encoding, and more specifically to a scalable rate control system for a video encoder.\n2. Description of the Related Art\nVideo standards are continuously being improved or otherwise replaced to improve performance and quality for a variety of applications. For example, the Advanced Video Coding (AVC) standard, Part 10 of MPEG4 (Motion Picture Experts Group), otherwise known as H.264, includes advanced compression techniques that were developed to enable transmission of video signals at a lower bit rate or to enable improved video quality at a given transmission rate. The H.264 standard is used for exemplary embodiments of the present invention although it is understood that the present invention is not limited to H.264 and is applicable to other video standards. The newer standard outperforms video compression techniques of prior standards in order to support higher quality streaming video at lower bit rates and to enable internet-based video and wireless applications and the like. The standard defines the syntax of the encoded video bit stream along with a method of decoding the bit stream. Each video frame is subdivided into one or more slices and encoded at the macroblock (MB) level, where each MB is a 16×16 block of pixels. The size of each slice is arbitrary and may range between a single MB up to all of the MBs in the frame. Each frame is also subdivided according to a rate control interval, specified as a number of MBs per interval, where the rate control interval also ranges between a single MB up to all of the MBs in the frame. The slice size and rate control interval size may be, but are not necessarily, the same.\nFurthermore, newer video and image codec standards are being developed to enable video systems having a wide range of computational capabilities to seamlessly process various frame rates, image resolutions, system bandwidth and bit rates in a heterogeneous environment. Scalable Video Coding (SVC), for example, is an extension of the H.264 main profile which addresses coding schemes for reliably delivery of video to diverse clients over heterogeneous networks using available system resources, particularly in scenarios where the downstream client capabilities, system resources, and network conditions are not known in advance, or dynamically changing from time to time. SVC achieves scalability by using a base- and an enhance-layer concept, where the enhanced layer is scalable from the base layer, and the rate control interval of base layer is scaled up too for spatial scalability. The base layer should be the simplest form in quality, spatial- and temporal-resolution. Whereas H.264 has relatively limited scalability, SVC provides multiple levels of scalability including temporal scalability, spatial scalability, complexity scalability and quality scalability. Temporal scalability generally refers to the number of frames per second (fps) of the video stream, such as 7.5 fps, 15 fps, 30 fps, etc. Spatial scalability refers to the resolution of each frame, such as common interface format (CIF) with 352 by 288 pixels per frame, or quarter CIF (QCIF) with 176 by 144 pixels per frame, although other spatial resolutions are contemplated, such as QVGA, VGA, SVGA, D1, HDTV, etc. Complexity scalability generally refers to the various computational capabilities and processing power of the devices processing the video information. Quality scalability is generally measured with a peak signal-to-noise (PSNR) metric defining the relative quality of a reconstructed image compared with an original image. Quality scalability is varied with the applicable bit rate, such as 64 kilobits per second (kbps), 128 kbps, 256 kbps, 512 kbps, etc.\nOne of the primary functional blocks within a video encoder for addressing scalability is the rate controller. It is desired to provide an efficient rate controller which is capable of generating a quality image in a heterogeneous and seamless environment in which devices with various processing capabilities demand scalable bit rates, image resolutions, and frame rates."}
{"text": "As a conventional technique, a configuration of a self discharge-type circuit closer is disclosed in which when the voltage between a pair of electrodes reaches a dielectric breakdown voltage in a circuit including a charged capacitor or the like, the electrodes are short-circuited to close the circuit (e.g., see Patent Document 1).\nAnother configuration is disclosed in which the contact elements of a pair of electrodes disposed in a vacuum vessel are brought into contact with each other to bring the electrodes into a conduction state (e.g., see Patent Document 2).\nFurther, as a circuit closer using a vacuum interrupter as with the aforementioned configuration, a configuration of a pulse-type circuit closer is disclosed in which: a pair of main electrodes is oppositely disposed in a vacuum vessel; a high voltage is applied to a trigger electrode provided in proximity to one or both of the main electrodes, to cause a micro-discharge between the trigger electrode and the electrode(s); and the plasma generated by the micro-discharge is injected into the main electrode(s) to break down the insulation between the main electrodes and to generate an arc, thus bringing the main electrodes into a conduction state (e.g., see Patent Document 3).\nIn addition, vacuum gap breakdown characteristics in the case of using oxygen-free copper as the electrode material are disclosed (Non-Patent Document 1, cited as FIG. 6 in the accompanying drawings of the present application). Dielectric breakdown voltages in the case of using different electrode materials are also disclosed (Non-Patent Document 2, cited as FIG. 7 in the accompanying drawings of the present application)."}
{"text": "In surgical procedures, sutures are commonly used to close incisions and to reunite damaged tissue. Typically, the sutures are maneuvered and passed through the affected tissue and the free ends of the sutures are individually tied together by the surgeon. In some surgical procedures, the surgical site area is sufficiently exposed to permit the surgeon to access and tie the suture manually with a surgical knot. In other surgical procedures, such as endoscopic procedures, laparoscopic procedures, arthroscopic procedures and the like, or when robotic surgical procedures occur, the surgical site is inaccessible to the surgeon's hands. As a result, the surgeon must tie each of the suture ends into a knot at a location remote from the surgical site, and then manipulate suitably configured instruments for sliding the surgical knot to the site of the incision. Further, surgeons may tie surgical knots intracorporeally (inside of the body) using surgical tools to tie the knot down to the tissue. In general, suture knot tying is cumbersome and is one of the more time-consuming steps in the suturing process of the surgical procedure. In the foregoing circumstances, it is desirable to replace knot tying during surgical procedures in order to significantly reduce the duration of surgical operations with a device or method that is simple for the surgeon to utilize. This is especially true with regard to minimally invasive surgical procedures where the tying of surgical knots within confined spaces is extremely difficult and time consuming.\nAdditionally, it is noted that knots create weak points in a suture. That is to say, when a failure load is applied to a knotted suture, assuming the suture is otherwise free from imperfections, the suture will break at the knot. Therefore, elimination of surgical knots in the suture would also eliminate the weak stress points created in the suture by the surgical knot.\nSuture locking and suture anchoring devices such as suture clips, surgical fasteners, hinged clips, suture terminating devices, hemostatic clips, and suture fixation devices of various configurations, designs, structures, and materials of construction are well known in the prior art. For example, U.S. Pat. No. 2,075,508 discloses a suture retainer whereby a suture may be fixed relative to a surgical button. The suture is received and wedged between the button and a clamping plate in order to securely clamp the suture. There are a number of shortcomings to this retainer. First, due to the suture manipulation required to use this retainer, a laparoscopic applicator device would be extremely difficult to produce profitably. Furthermore, the wedging and clamping action will induce stress concentrations in the suture, which likely lead to reduced failure loads.\nU.S. Pat. No. 6,066,160 to Colvin, et al. discloses a suture terminator device for use in minimally invasive surgery. The suture terminator device includes a pair of locking apertures with teeth for engaging a portion of a suture at the locking apertures' threaded ends. Not only does this device require laparoscopically guiding two strands of suture between two very small apertures, it also requires that the sutures be squeezed tightly by these teethed apertures, thereby diminishing the integrity of the suture and significantly reducing the maximum tension the suture can withstand.\nU.S. Pat. No. 6,106,545 to Egan discloses a suture tensioning and fixation device, which includes a retaining element for frictionally engaging a suture that may subsequently be melted to bond to the suture for a permanent fixation. This melting and bonding action will compromise the integrity of the suture and therefore, because suture strength is of utmost concern in most surgeries, this method of fixation of the retaining element to the suture is not suitable for most surgeries.\nU.S. Pat. No. 5,474,572 to Hayhurst and U.S. Pat. No. 5,645,553 to Kolesa et al. disclose the use of a hinged clip that snaps closed after the suture threads are placed within the holding members. The hinge clip is snapped into place such that the suture is held transversely across the holding members, thus locking the suture in place. There is a possibility of improperly actuating or inadvertently releasing the snap, which could lead to an insecure fixation of the suture. In addition, weak stress points are created where the suture is held within the clip.\nIn view of the deficiencies of the prior art discussed hereinabove, there remains a need for a suture anchoring device that is simple to use, particularly during laparoscopic surgery, in order to eliminate manual knot tying by the surgeon performing the surgical procedure, while not compromising the integrity of the suture.\nAnother object of the present invention is to provide a suture anchoring device that is at least as strong as conventional surgical knots.\nAnother object of the present invention is to provide a suture anchoring device that is easy and cost effective to manufacture in bulk.\nAnother object of the present invention is provide a suture anchoring device that is suitable for a wide variety of sutures (i.e., such as monofilament and braided sutures)."}
{"text": "1. Field of the Invention\nThe present invention relates to products made from recycled materials, in particular to a process for making stronger products from recycled materials in a municipal waste stream specifically using a venturi extrusion methodology,\n2. Description of the Related Art\nMunicipal waste has become a serious problem in western society. Large fleets of trucks are used to transport the waste products of modern living to landfills on the edge of every modern city. The requirement for landfills is eating up increasingly large sections of land that could be used for other purposes. There is also the tremendous amount of recyclable material that is not being recycled. There is a very critical need to utilize this vast resource and at the same time save the “waste” land that is now being used as landfills.\nThere have been numerous attempts to create useful products from recycled materials. The main components in the majority of them have been waste wood and plastics. The methodologies utilized include injection molding and extrusion.\nIn current processes for making products from recycled materials the main problem has been the lack of strength of the final product. If the products are to be used for such applications as railroad ties, where immense pressures are applied, then the strength of the final product must be very high. However with the current methodologies most recycled building products lack sufficient tensile strength to be used for load bearing purposes.\nRecycled products that have not been strengthened usually can only be used for decorative purposes due to lack of tensile strength. The exceptions are typically recycled wood based products. Without sufficient tensile strength the products will shear or bend easily. Therefore most current recycled products cannot be uses for construction beams, unsupported spans, railroad ties, etc, where greater strength is necessary.\nWhat is needed is some method for improving the strength of recycled products that solves one or more of the problems described herein."}
{"text": "Switching Mode Power Supplies (SMPS) are used in a variety of portable electronic devices including laptop computers, cellular phones, personal digital assistants, video games, video cameras, etc. They may convert a dc signal at one voltage level to a dc signal at a different voltage level (this is a dc-dc converter), an Alternating Current (ac) signal to a dc signal (this is a an ac-dc converter), a dc signal to an ac signal (this is a dc-ac converter), or an ac signal to an ac signal (this is an ac-ac converter). Generally, switching mode power supplies have included switching power supply controllers or converters that operate in a continuous mode under heavy load conditions, i.e., a load that draws a large current, and under a skip mode or pulse skip mode under light loading conditions, i.e., a load that draws a small current. In the past, semiconductor manufacturers used various methods and structures to form the switching power supply controllers, such as Pulse Width Modulated (PWM) power supply controllers that regulated the value of a voltage supplied by a power supply system. In some cases, the switching power supply controllers were capable of operating in a fixed frequency or continuous operating mode during normal operation. When the current drawn by the load that was receiving power from the power supply system decreased, some of the prior switching power supply controllers operated in a light load mode that skipped some of the PWM cycles. The light load operating mode has been referred to as a skip mode or a burst mode. When the load again required a higher current, the switching regulator circuit exited the skip mode and returned to normal operation.\nUsually the converters include a compensation network that is connected to an error amplifier to stabilize the system and to optimize the transient response based on a small signal behavior model. However, when a large signal transient event occurs such as a step load transient from a light load condition to a full load condition, the converter is unable to achieve the desired response due to saturation recovery of the compensation network and the slew rate limitation of an error amplifier.\nFIG. 1 is a prior art schematic diagram of a portion of a power supply system 10 that includes a switching power supply controller 12. System 10 receives power between a power input terminal 14 and a power return terminal 16 and forms an output voltage (VOUT) between output 18 and terminal 16. Controller 12 is configured to regulate output voltage VOUT to a desired value or target value within a range of values around the target value. For example, the target value may be five volts (5 v) and the range of values may be plus or minus five percent (5%) around the five volts. System 10 usually includes a power switch such as a power transistor 20 and a rectifier 22 that are connected to control an inductor current I24 that flows through an inductor 24. Rectifier 22 may be a synchronous Metal Oxide Semiconductor Field Effect Transistor, a diode, or the like. A capacitor 26 is connected between output 18 and terminal 16 in order to assist in forming output voltage VOUT. A voltage sense network 28 may be coupled to output 18 to provide a voltage sense signal VS at node 30 that is representative of the instantaneous value of output voltage VOUT. By way of example, voltage sense network 28 is comprised of resistors 32 and 34 having terminals that are commonly connected together to form node 30. In addition resistor 32 has a terminal connected to output 18 and resistor 34 has a terminal connected to power return terminal 16. Voltage sensing network 28 may be any type of sensing network that provides sense signal VS at node 30 that is representative of the value of output voltage VOUT. A load 36 is generally connected between output 18 and terminal 16 in order to receive output voltage VOUT and to receive a load current ILOAD. It should be noted that load current ILOAD is the sum of current I24 and a current I26 that may flow from capacitor 26.\nSwitching power supply controller 12 receives operating power from a regulator 35 that is connected between a voltage input 38 and a voltage return 40. Input 38 and return 40 typically are connected to respective terminals 14 and 16. It should be noted that regulator 40 may provide reference voltage VREF. Controller 12 is configured to form a switching drive signal on an output 42 that is suitable for driving and operating transistor 20 to regulate the value of output voltage VOUT. Voltage sense signal VS from voltage sense network 28 is received by controller 12 on a voltage sense input 44.\nController 12 includes a PWM control module 50 suitable for generating a PWM switching signal that is input into a buffer driver or buffer 52. Buffer 52 has an output terminal connected to a gate terminal of power transistor 20. Controller 12 further includes a feedback network 54 that comprises an operational amplifier 56 and a compensation network 58. By way of example, compensation network 58 is a passive voltage compensation network. More particularly, operational amplifier 56 serves as an error amplifier that has an inverting input terminal, a non-inverting input terminal, and an output terminal, where the non-inverting input terminal is coupled for receiving a reference voltage VREF, the inverting input terminal is coupled to its output terminal and to voltage sense node 44 through compensation network 58. By way of example, compensation network 58 is composed of a resistor 60 connected between the inverting input terminal of operational amplifier 56 and voltage sense node 44, and a resistor capacitor network 62 coupled between the inverting input terminal and the output terminal of operational amplifier 56. Resistor capacitor network 62 is comprised of a capacitor 64 coupled in parallel with a resistor 66 and a capacitor 68 which are connected in series. The output terminal of operational amplifier 56 is directly connected to an input terminal of PWM control module 50.\nIn operation, power supply system 10 typically operates in one of two operating modes: a continuous operating mode or a pulse skip (or burst) operating mode. Under a heavy or non-light load condition, PWM control module 50 operates at its nominal or full operating frequency and inductor current I24 is continuous. Under a light load or a no load condition, load current ILOAD decreases and inductor current I24 becomes discontinuous. If the pulse skip mode is enabled, the operating frequency or switching frequency at the output terminal of PWM control module 50 is reduced in response to the decrease in the loading current, thereby reducing power dissipation.\nFIGS. 2a, 2b, 2c, and 2d are plots that illustrate various signals that are generated by controller 12 when it operates in a continuous operating mode. The abscissas of plots 2a, 2b, 2c, and 2d indicate time and the ordinates of plots 2a, 2b, and 2c indicate voltage, whereas the ordinate of plot 2d indicates current. More particularly, plot 2a illustrates the voltage VCOMP transmitted from the output terminal of operational amplifier 56 to the input terminal of PWM control module 50; plot 2b illustrates output voltage VOUT that appears between output 18 and terminal 16; plot 2c illustrates the voltage VSWN appearing at node 25; and plot 2d illustrates inductor current I24. In FIG. 2, controller 12 operates in the continuously pulsing PWM mode, thus inductor current I24 is continuous. Under this condition, operational amplifier 56 does not operate in saturation and there is little change in the DC bias of the capacitors of resistor capacitor network 62, i.e., capacitors 64 and 68. More particularly, in the steady state continuously pulsing mode of operation the DC bias across capacitor 68 is substantially equal to the difference between the average voltage level at the output terminal of error amplifier 56 during the continuously pulsing mode of operation and reference voltage VREF that appears at the non-inverting input terminal of error amplifier 56, i.e., VC68=VCOMP—AVG−VREF. In the pulse skipping mode, when output voltage VOUT is higher than reference voltage VREF, voltage VCOMP remains at its minimum level. Under this condition, the DC bias voltage across capacitor 68 is substantially equal to the difference between the minimum voltage value of voltage VCOMP and reference voltage VREF, i.e., VC68=VCOMP—MIN−VREF.\nWhen there is a step-up load transience, voltage VCOMP increases to a value that is substantially equal to the average voltage level at the output terminal of error amplifier 56. Because the DC bias voltage across capacitor 68 cannot change instantaneously, a voltage difference substantially equal to VCOMP−VC68−VREF is added across resistor 66, which causes a droop current to be injected into node 59 and through resistor 60. Voltage VC68 is the voltage across capacitor 68. The additional droop current results in an additional voltage droop and a longer time for output voltage VOUT to recover.\nFIGS. 3a, 3b, 3c, and 3d are plots that illustrate various signals that are generated by controller 12 when there is a step-up load transience, the skip mode is enabled, and controller 12 has been operating in the skip mode. The abscissas of plots 3a, 3b, 3c, and 3d indicate time and the ordinates of plots 3a, 3b, and 3c indicate voltage, whereas the ordinate of plot 3d indicates current. More particularly, plot 3a illustrates the voltage VCOMP transmitted from the output terminal of operational amplifier 56 to the input terminal of PWM control module 50; plot 3b illustrates the output voltage VOUT appearing between output 18 and terminal 16; plot 3c illustrates the voltage VSWN appearing at node 25; and plot 3d illustrates inductor current I24 and load current ILOAD.\nBecause controller 12 went into the skip mode of operation, the system transient response is degraded due to the DC level deviation in feedback network 54.\nAccordingly, it would be advantageous to have a method and a power supply controller with a fast transient response under heavy and light loading conditions. It would be of further advantage for the circuit and method to be cost efficient to implement."}
{"text": "Various kinds of inhaler are known for delivering metered doses of medicament for inhalation. Typically, pressurized liquid systems and dry powder systems are used to deliver medicament to a patient. In one pressurized liquid system a liquid containing medicament is stored in a pressurized container and a metered dose is delivered as an aerosol by a pressurized metered dose inhaler. In dry powder systems powder containing medicament is stored in individual blisters, on tape or in bulk in a chamber and a metered dose is provided by a dry powder inhaler.\nAll such inhalers include parts which come into contact with medicament and parts, some in common with the parts which come into contact with medicament, which come into contact with the mouth or nose of a user. It will be appreciated that it is important to ensure that those parts do not have any adverse effect on either the medicament or the user.\nIt is an aim of the present invention to provide an improved inhaler and a method of manufacturing a housing of the same in view of this requirement."}
{"text": "Poly(vinyl chloride) (“PVC”) is commonly utilized as a material of construction for consumer and industrial goods. PVC possesses the advantageous features of low cost, durability, moisture resistance, tailored stiffness, dimensional stability, and flame retardancy. Virgin PVC is readily processed into sheet, tubes, and other forms using conventional processing equipment such as extruders and thermal compression bonding equipment. The reported density of poly(vinyl chloride) is 1.45 g/cm3.\nPlasticizers can be added to PVC, which is known to be rigid in the absence of such additives, to make it more flexible and more suitable for an even broader range of applications, including but not limited to applications such as plumbing and electrical cable insulation. The addition of plasticizers and other such flexibilizing agents will lower the modulus and the density of the PVC.\nThe use of chemical and physical inert gas blowing agents can also be used to lower the density of PVC substrates.\nWhile PVC is attractive for many first life commercial uses, the opportunities to reclaim and reuse PVC for subsequent future life applications are limited. The recycling of reclaimed PVC possesses inherent challenges due both to the processing difficulties caused by the various and often unknown additives and fillers that may have been used during the first life and to the sorting and separation of reclaimed PVC to generate a more homogenous raw material for use in future applications. Specifically, a significant amount of reclaimed PVC sources incorporate mixtures of other materials or components, for example pigments, colorants, fillers, plasticizers and the like, that limit PVC's potential reuse in future applications. Reclaimed PVC materials introduce feedstock variability for which conventional melt processing may be ill-suited and which limits the commercial utility for future life applications of PVC. This challenge becomes more acute as one tries to use a high fraction of reclaimed PVC within conventional processing methods."}
{"text": "Disposable capsules are known since time of the type comprising a container having an plane openable or punchable base for the introduction of water and a closed plane base from a plane filter, protected by a peelable film before use, through which the beverage produced by infusion of water with the aromatic substance present in the container is delivered.\nIt is known that such capsules, due to their conformation, can negatively affect the correct execution of the infusion, with a consequent penalization of the organoleptic characteristics of the produced beverage.\nThis is often due to the characteristics of the filter, having a surface extension sometimes insufficient, being substantially equal to that of the plane base of the container to which it is applied.\nWith particular reference to capsules for American coffee, which require an extraction procedure of the aroma purely by percolation, the filter, if having an insufficient surface, tends to obstruct and form a non desired overpressure inside the container. The provision of a greater filtering surface in a known disposable capsule would require a more bulky container in a transversal direction, which in turn would require specifically dimensioning also the infusion chamber in which the capsule is housed."}
{"text": "1. Field of the Invention\nThis invention relates generally to computer-implemented systems for optimal query execution in relational database processing systems and particularly to a system for optimizing query ordering requirements by normalization of relational order specifications to accommodate predicates and functional dependencies before order detection in a query compiler system.\n2. Description of the Related Art\nA relational database management system (RDBMS) is a computer-implemented database processing system that uses relational techniques for storing and retrieving data. Relational database management systems are computerized information storage and retrieval systems in which data in the form of tables (formally denominated \"relations\") are typically stored for use on disk drives or similar mass data stores. A \"relation\" includes a set of rows (formally denominated \"tuples\" or \"records\") spanning several columns (formally denominated \"attributes\"). Each \"attribute\" in a relation includes \"restrictions\" on the data contents thereof, which may include designation as a primary or foreign key. A \"tuple\" expresses a mathematical relation between its \"attributes\" (column elements), while a \"record\" does not. Reference is made to C. J. Date, An Introduction to Database Systems, 6th Edition, Addison-Wesley Publishing Company, Reading, Mass. (1994) for a general treatment of the relational database art.\nA relational database processing system is structured to accept commands to store, retrieve and delete data using high-level query languages such as the Structured Query Language (SQL). As used herein, a \"query\" refers to a set of user commands for retrieving data from a stored database. The query language requires the return of a particular data set (relation) in response to a particular query but the method of query execution employed by the database management system is not specified by the query. The method of query execution is herein denominated the \"Query Execution Plan (QEP)\" and there are typically many different useful QEPs for any particular user query, each of which returns the required data set. For large databases, the execution plan selected to execute a query must provide the required data return at a reasonable cost in time and hardware resources. Most RDBPSs include a query compiler to translate queries into a QEP, which includes a query optimizer to select an efficiently executable QEP.\nAmong many well-known SQL commands, a SQL query may include the ORDER BY clause. For instance, the following SQL query example returns a list of movie titles, in alphabetical order, from a relation keyed to movie stars:\nSELECT DISTINCT MOVIE.sub.-- TITLE\nFROM MOVIE.sub.-- STARS PA1 ORDER BY MOVIE.sub.-- TITLE ASC;\nIn this example, only one attribute or column is used for ordering. ASC requires an ascending order, as opposed to DESC for descending order. This query imposes an \"order requirement\" that must be satisfied by the query compiler when it generates a QEP. Moreover, the query compiler itself may generate internal order requirements arising from combinations of other SQL operations. For example, the internal procedures used to satisfy GROUP BY and DISTINCT clauses generally require an input record stream with an \"order property\". The sort-merge join operation is known in the art to require an \"order property\" in each of its record stream inputs. Also, where VIEW specifications are used, an order requirement may arise for the intermediate results of a query as well as for the final query data return. As used herein, the terms \"order property\" and \"order requirement\" refer to two order specifications imposed on a group of records in a relation. Such order specifications list one or more \"attributes\" (columns) within each record on which the records are ordered.\nA typical RDBPS puts records in order by either sorting the records on their \"order requirement\" attributes or, when possible, by scanning an existing index that already provides the required ordering. Sorting is a costly operation that is avoided when possible, so the typical query compiler first searches for an existing index that satisfies the order requirement. As used herein, an \"index\" is a user-defined relation that maintains a certain order of records. Each relation may be associated with any number of indexes or with no index. When no index can be found, sorting is unavoidable and the query compiler then attempts to minimize the number of sort \"attributes\" to optimize execution plan efficiency. Thus, as it generates an execution plan, a query compiler must detect whether a sort can be avoided and, if unavoidable, how many attributes in the \"sort requirement\" may be eliminated during actual sorting. This process is herein denominated \"order detection\". The following example illustrates the \"order detection\" problem known in the art:\nSELECT SALARY, NAME\nFROM EMP\nWHERE SALARY=50,000\nAND DEPTNO=:HV\nORDER BY 1,2;\nIn this SQL example, each execution plan must produce a stream of records in the order specified by the ORDER BY clause. Consequently, the query compiler must satisfy the order requirement (SALARY, NAME), with SALARY being the major ordering attribute. If an ascending index exists on attributes (SALARY, NAME), then an index scan produces the required \"order property\" without sorting. Alternatively, an unordered table scan followed by a sorting process may be employed to produce the required \"order property\" in the final relation. Typically, a query compiler considers both plans and chooses the less expensive to execute.\nIn the above example, if an ascending index also exists on the NAME attribute, then such index could also be used to produce the requisite order property without sorting. No sorting is actually then required because the SALARY attribute is bound to a constant by the predicate SALARY=50,000. Accordingly, the SALARY attribute does not affect the order property because all qualifying records have the same SALARY attribute (the same value in the SALARY column). Similarly, where an unordered table scan and sort are used, the sort need only be on the NAME attribute. Finally, if an index is found for (DEPTNO, NAME), it also satisfies the order requirement without sorting because, like the SALARY attribute, the DEPTNO attribute is bound to a constant (in this case to a host variable value).\nThe above-described example shows how the order detection process could improve execution efficiency if it could be aware of which relation attributes are \"bound\" to constants by way of predicates during execution plan generation. Order detection could also improve efficiency if it could be aware of keys and any functional dependencies arising therefrom, as illustrated by the following example:\nSELECT E.EMPNO, E.NAME, D.DEPTNO\nFROM EMP E, DEPT D\nWHERE E.DEPTNO=D.DEPTNO\nORDER BY 1,2,3;\nIn this SQL query example, the query compiler must satisfy the order requirement vector (E.EMPNO, E.NAME, D.DEPTNO), which includes the first, second and third attributes in each record. If the EMPNO attribute is a key of the EMP relation, then sorting on EMPNO alone is sufficient to meet this order requirement. Alternatively, using an index ordered on EMPNO is sufficient without sorting. The skilled database practitioner can appreciate the reason for this sufficiency, which arises because the primary key EMPNO creates a \"functional dependency\" between the EMPNO attribute and all other attributes within the relation. This functional dependency combined with the predicate E.DEPTNO=D. DEPTNO means that E.EMPNO \"functionally determines\" both E.NAME and D.DEPTNO, which in turn implies that sorting on E.EMPNO is sufficient to create the three-attribute order property specified by the above example. Until now, the sort detection processes in the art have not been aware of these relationships.\nSuccinctly, a \"functional dependency\" herein denotes a relationship between two attributes where \"A functionally determines B\" when every record in a relation that has a particular value for attribute A has an unchanging value for attribute B. In the art, this functional dependency is denoted by the standard notation A.fwdarw.B. An attribute B is \"functionally dependent\" on a set of attributes or columns C={C.sub.1, C.sub.2, . . . , C.sub.n } if, for any two records with the same values for the columns C, the value of column B is the same. Equivalently, it is said B is functionally determined by C, which is denoted as C.fwdarw.B. As is known in the art, key attributes are actually a special case of functional dependency referred to as a \"key dependency\". A key dependency functionally determines all columns of the record stream for which it is a key. Any remaining dependencies are denoted \"non-key dependencies\". Functional dependencies usually arise from keys, whether primary or foreign. Joins and other relational operations can transform a key dependency into a non-key dependency. Relational operations such as DISTINCT and GROUP BY can add key dependencies to derived relations. For the purposes of this disclosure, it is assumed that means are available to determine whether C.fwdarw.B is true.\nReference is made to Hugh Darwen et al. (RELATIONAL DATABASE WRITINGS 1989-1991, C. J. Date with Hugh Darwen, Chapter 10: The Role of Functional Dependencies in Query Decomposition for a description of the semantic concept of functional dependencies. Darwen et al. suggest algorithms for deriving functional dependencies from base relations and for propagating functional dependencies across derived relations. Darwen et al. also consider a variety of uses for functional dependencies by a RDBMS query compiler, but they neither consider nor suggest how to accommodate functional dependencies during order detection.\nAs is known in the art, a query compiler for a relational database processing system translates a high-level query into a query execution plan (QEP), which is then interpreted by a database engine at run time. Conceptually, a QEP is viewed in the art as a data flow graph of operators, whese each node in the graph corresponds to a set of output records arising from a relational operation such as \"join\" or a lower-level operation such as \"sort\". FIG. 2 shows a query execution plan (QEP) for the following exemplary SQL qudry.\nSELECT NAME, ADDRESS\nFROM EMP E, DEPT D\nWHERE E.DNO=D.DNO AND MGR=`Haas`\nNote that the QEP in FIG. 2 is a directed graph of LOw LEvel Plan OPerators (LOLEPOPs), herein formally denominated \"relation nodes\". Note that the data-flow arcs exemplified by arc 12 point toward the source of the record stream and not in the direction of data-flow. For instance, the SORT node 14, representing a relational SORT operation, receives a record stream from the ACCESS node 16 along the directed data-flow arc 18 and produces a directed record stream along arc 12 to the JOIN node 20. For the purposes of this disclosure, notice that each relation node in FIG. 2 is coupled to one or more other relation nodes by data-flow arcs (herein formally denominated \"directed record streams\") exemplified by directed record streams 12 and 18. With the notation used herein, each relation node N.sub.i consumes one or more incoming record streams {R.sub.ji } and produces at least one output record stream R.sub.ik that is directed to one or more other nodes (k and j may each assume more than one value). In the example shown in FIG. 2, each relation node produces only one output record stream routed as an input record stream to one other relation node.\nReference is made to Guy M. Lohman (\"Grammar-like Functional Rules for Representing Query Optimization Alternatives\", A CM SIGMOD International Conference on Management of Data, 1988, pp. 18-27) for a detailed discussion of the \"properties\" associated with relation nodes or operators in a QEP. Every relation node is associated with a set of properties that summarizes the work accomplished in the relation thus far in the dataflow path from bottom (base-tables) to top (TOP node). Lohman et al. discuss these properties in terms of the query execution cost model and notes that all properties are handled uniformly as elements of a \"property vector\", which is easily extended to include more properties. The \"property vector\" includes the \"attributes property\", herein denominating the columns of a record stream, which may include columns from base-relations and those derived from expressions. The \"attributes property\" is used to track the record stream columns.\nExamples of properties include the set of attributes that make up each record or tuple in the stream, represented by the head of each node. Other examples include the set of predicates that are applied to the incoming records within the node body, the specification of which attributes form unique keys and/or functional dependencies (if any), the order specification for the record stream, and the like. Each relational operator node in a QEP determines the properties of its output record stream.\nThe properties of an output record stream from a node are a function of the node input streams and the internal operations applied within the node. For example, a sort operation passes on all of the properties of its input stream unchanged except for the order property. A query compiler typically constructs a QEP from the bottom up, operator by operator, computing and updating properties as it proceeds. Depending on the implementation, some properties may be stored in QEP operators (within \"nodes\") while others may be recomputed when needed to save space.\nEach node in a QEP receives one or more directed record streams each having an \"order property\". By virtue of operations within the node, each node produces an output record stream according to an \"order requirement\". As discussed above, the ORDER BY, GROUP BY, and DISTINCT specifications create order requirements. Moreover, the query compiler may generate order requirements, depending on the alternative execution plan characteristics. An example of an execution strategy is the \"sort merge join\" operation, which requires certain order properties in the incoming record streams. There may be multiple order requirements imposed within any given SQL query. Both order properties and order requirements are represented by an order specification and are herein referred to generally as order specifications. As a simple example, ORDER BY C.sub.1, C.sub.2, C.sub.3 gives rise to the order requirement vector (C.sub.1, C.sub.2, C.sub.3). In the case of the GROUP BY and DISTINCT clauses, any permutation of the order specification suffices. For example, GROUP BY C.sub.1, C.sub.2 is represented by either of the order requirements (C.sub.1, C.sub.2) and (C.sub.2, C.sub.1). The order requirement of the GROUP BY operation is satisfied if either order requirement is satisfied. In some query compilers, these two order requirement vectors can be represented as a single order requirement specification. For the purposes of this disclosure and without loss of generality, it is assumed that each record stream must satisfy only a single order requirement. In the general case, many order requirements may be imposed on a single record stream, which might benefit from additional procedures to \"compact\" multiple order requirements into a single order specification.\nAs a QEP is constructed, order requirements must be satisfied. That is, for instance, when an ORDER BY operator is added to the QEP, the \"order detection\" procedure examines the \"order property\" of the vector associated with the input stream to determine whether it satisfies the \"order requirement\" imposed by relational operations in the local relation node. If the \"order property\" of the incoming record stream does not satisfy the \"order requirement\" imposed by the receiving relation node, then the query compiler adds a sort operator to the QEP, representing an additional operation internal to the receiving relation node. This sort operator adds substantial cost to the resulting QEP and it is a desired function of the \"order detection\" procedure to avoid adding sort operators wherever possible.\nInitially, the properties of stored base tables and access methods are determined from the system catalogs. At first, no predicates are applied and no cost has been incurred in the query. The \"order property\" is unknown unless the base relation is known to store records in some order or if an index is used to access the relation in which case the order is defined by the set of attributes on which the records are ordered.\nDuring compilation, the predicates in a SQL query are usually broken down into conjunctive normal form. This permits each conjunct to be applied as an independent predicate. When a predicate or conjunct is applied to an input record stream, each record in the resulting output record stream is imprinted with the \"property\" that it satisfies the predicate. Herein, the term \"predicate property\" denominates the property used to track all predicates that have been applied to a record stream.\nThe \"order property\" for a record stream represents the physical ordering of the stream's records. An order property is denoted by listing the attributes or columns of the order in major to minor sequence. Such representation is herein denominated an \"order specification vector\". For example, (C.sub.1, C.sub.2, C.sub.3) is an \"order specification vector\" with C.sub.1 as the most major ordering column and C.sub.3 is the most minor ordering column. Column C.sub.2 is in a \"more major\" position relative to C.sub.3 and in a \"more minor\" position relative to C.sub.1. Two records are ordered with respect to an \"order requirement\" vector by comparing attribute values in most major to most minor sequence until order is determined (ISO-ANSI, \"ISO-ANSI Working Draft: Database Language SQL2 and SQL3; X3H2; ISO/IEC/JTC1/SC21/WG3, 1993\"). A record stream can acquire an \"order property\" because of an index access or a sort operation while other operations may preserve the order of an input record stream. For example, a nested-loop join operation preserves the order of the outer relation.\nThe \"functional dependency\" property specifies a collection of functional dependencies that hold for the attributes of the record stream. For base relations, these are typically the key dependencies of the base relations. Operations that remove duplicates can add key dependencies as described above. Key dependencies survive joins as non-key functional dependencies. The above-cited Darwen et al. reference describes the functional dependency derivation and propagation for relational operations. For the order detection process, the \"functional dependency\" property permits enumeration of all functional dependencies holding for the record stream.\nThere is an untilled need in the art for a sort detection process that permits the query compiler to understand and exploit the effects of both predicates and functional dependencies during order detection. The failure in the art to account for these effects causes the query compiler to inject unnecessary sorts or to sort on more attributes than necessary. This in turn causes the compiler to generate suboptimal query execution plans. The desire to reduce QEP cost in large database systems motivates the need for any useful method that avoids unnecessary sort operations during execution plan generation.\nPractitioners in the relational database art generally appreciate the absence of a solution to this problem. For instance, Boykin et al. (\"Order Class Detection\", IBM Technical Disclosure Bulletin, Vol. 32, no. 52, pp. 283-285, October, 1989) discuss the continuing need to improve order detection during query optimization. Boykin et al. propose a process for recognizing when special combined order classes can be constructed to satisfy the order requirements arising from more than one order specification in a query. An order class is a list of attributes or columns, each having the attribute \"ascending\", \"descending\", or \"either\". The combined order classes are used to locate indexes or to specify the sorting operation necessary to satisfy the order requirements. As used herein, the \"order class\" is denominated an \"order specification vector\", which may be an \"order property vector O.sub.Pi \" or an \"order requirement vector O.sub.Ri \". By combining the \"order specification vectors\" within a query, Boykin et al. permit the query compiler to insert a single sort to satisfy two separate ordering requirements, but do not consider the effects of functional dependencies or predicates on sort specifications.\nL. W. Yongren (\"Dynamic Ordering of Joined Rows using Fields from Multiple Tables\", IBM Technical Disclosure Bulletin, Vol. 36, no. 11, pp. 363-366, November 1993) disclose a method for ordering joined \"live\" data (dynamic views) without producing a temporary result table. Yongren proposes maintaining indexes that can be updated responsive to updates to the underlying base-relations. By providing for a \"composite index\", Yongren is able to avoid resorting operations during the maintenance of dynamic VIEW tables, but he neither considers nor suggests methods for eliminating sort operations or index scans by exploiting functional dependencies.\nSimilarly, in U.S. Pat. No. 5,241,648, Cheng et al. propose a hybrid method for joining tables that relies on maintaining indexes to speed the nested join procedure. In U.S. Pat. No. 5,369,761, Conley et al. speed the table join operation by selectively \"denormalizing\" the base-relations to avoid joining unnecessary attributes. Neither Conley et al. nor Cheng et al. consider or suggest methods for eliminating sort operations by detecting redundancy arising from functional dependencies.\nThere is a need in the art for an order detection process that fully assimilates and exploits the effects of both predicates and functional dependencies to avoid sorting whenever possible. The absence in the art of any process for exploiting predicate and functional dependency effects during query order detection reduces query execution plan efficiencies. These unresolved problems and deficiencies are clearly felt in the art and are solved by this invention in the manner described below."}
{"text": "1. Field\nThe present disclosure relates to computer systems and methods in which data resources are shared among data consumers while preserving data integrity and consistency relative to each consumer. More particularly, the disclosure concerns implementations of mutual exclusion mechanisms such as reader-writer locking.\n2. Description of the Prior Art\nBy way of background, reader-writer synchronization is a mutual exclusion technique that is suitable for use in shared memory multiprocessor computing environments to protect a set of shared data. One type of reader-writer synchronization, known as reader-writer locking, allows read operations (readers) to share lock access in order to facilitate parallel data reads, but requires write operations (writers) to obtain exclusive lock access for writing the data. The technique is well suited to shared memory multiprocessor computing environments in which the number of readers accessing a shared data set is large in comparison to the number of writers, and wherein the overhead cost of requiring serialized lock acquisition for readers would be high. For example, a network routing table that is updated at most once every few minutes but searched many thousands of times per second is a case where serialized read-side locking would be quite burdensome.\nReader-writer locks are conventionally implemented using a single global lock that is shared among processors. This approach requires readers and writers to contend for one global lock on an equal footing, but produces memory contention delays due to cache line bouncing of the lock between each processor's cache. Insofar as reader-writer locks are premised on the existence of a read-intensive processing environment, readers may be unduly penalized, especially if their critical sections are short and their lock acquisition frequency is high. A distributed reader-writer lock approach is presented in Hsieh and Weihl, “Scalable Reader/Writer Locks for Parallel Systems”, 1991. It requires the readers to acquire only a local per-processor reader/writer lock that will usually reside in the memory cache of the processor that hosts the acquiring reader. However, the writers must acquire all of the local reader/writer locks, which degrades writer performance due to memory contention, and in some cases due to new readers being allowed to starve a writer while the latter is waiting for one of the local reader/writer locks. A further disadvantage associated with both non-distributed and distributed reader-writer locking is that lock acquisition imposes a burden on readers, even in the absence of a writer. Reader-writer locks are typically implemented as semaphores, mutex locks and spinlocks. Acquiring each of these lock types often imposes the cost of atomic instructions and/or memory barriers. In a read-mostly computing environment, the overhead associated with these operations falls mostly on readers.\nImproved read-side performance is provided by the locking technique disclosed in commonly-owned U.S. Pat. No. 7,934,062, which requires no read-side lock acquisition except when a writer announces its intention to acquire the reader-writer lock. However, the write-side performance of this method can be degraded in systems with many processors. This is because writers must wait for a grace period to elapse before acquiring the reader-writer lock. All processors must pass through a quiescent state that guarantees each reader will have an opportunity to note the writer's locking attempt, and thereby synchronize on the reader-writer lock.\nThe present disclosure introduces techniques for reducing writer latency in large multiprocessor systems that employ data synchronization mechanisms, such as the grace period-based reader-writer locking approach disclosed in U.S. Pat. No. 7,934,062 or the distributed locking scheme proposed by Hsieh and Weihl. A technique for reducing writer latency in a multithreaded user-mode embodiment of the Hsieh and Weihl distributed locking method is also disclosed. The techniques disclosed herein are also useful for other synchronization operations, such as expedited grace period detection in multiprocessor systems implementing read-copy update (RCU) synchronization."}
{"text": "This invention is directed to a wrench of the locking type which will grip objects of various sizes and hold these objects without the necessity of the wrench operator constantly applying force to the handles of the wrench. An over-center toggle type linkage is used to hold a lower pivoted jaw relative to an upper fixed jaw.\nThe primary object of this invention is to provide a wrench that automatically self-adjusts itself to various size objects. The wrench operator does not have to adjust a screw or other adjustment means before gripping an object.\nAnother object of this invention is to provide a wrench that in addition to self-adjusting to various sizes is also capable of gripping objects with a wide variation in clamping force. The degree of force applied to an object is controlled each time the wrench is closed or used by merely applying pressure on a special force control lever provided on the wrench.\nYet another object of this invention is to provide a wrench that lends itself to one hand operation. In one variation the wrench is provided with a force control lever on the top forward part of the upper handle, where it can be operated with ones thumb, while the wrench handles are closed by using ones hand and fingers. In standard manually adjusted wrenches of this type (A well-known brand or trade name of these standard wrenches is Vice-Grip.), the adjustment screw must be adjusted and often re-adjusted prior to use, and/or the adjustment screw must be turned with one hand while some force is applied to the wrench handles with the other hand to thereby arrive at a reasonable adjustment. In many circumstances (such as when parts must be held) this type of manual adjustment is very difficult or impossible."}
{"text": "The present disclosure relates to a light-emitting device and a method of manufacturing the same.\nIn a laser light source, realizing high power is a major task. Therefore, semiconductor optical amplifiers (SOA) have been studied as devices that amplify light from a laser light source as well as realizing high power of semiconductor laser devices. Here, an optical amplifier is a device that directly amplifies light without converting an optical signal into an electrical signal. The optical amplifier has a laser structure having no resonator and amplifies incident light with an optical gain of the amplifier.\nSince optical amplifiers have been developed mainly for optical communication, semiconductor optical amplifiers of a 405 nm band have rarely been put to a practical use. For example, Japanese Unexamined Patent Application Publication No. 5-067845 discloses a semiconductor optical amplifier of a 1.5 μm band which has a tapered ridge stripe configuration in which a GaInAsP-based compound semiconductor is used. According to the technology disclosed in Japanese Unexamined Patent Application Publication No. 5-067845, in the semiconductor optical amplifier, a mode field is expanded along an optical waveguide width by gradually widening the optical waveguide width in a tapered shape from an optical waveguide serving as a narrow incident side, which satisfies single-mode conditions, to an optical waveguide serving as an output side. Thus, the highest output of the semiconductor optical amplifier is designed to be improved.\nIn general, semiconductor laser devices or semiconductor optical amplifiers have a layer (In the field of laser diode, the word “layer” is used instead of lamination) structure formed on a base substrate. Here, the layer structure includes n-type compound semiconductor layers (specifically, an n-type clad layer and an n-type light guiding layer) doped with n-type impurities, an active layer, and p-type compound semiconductor layers (specifically, a p-type light guiding layer and a p-type clad layer) doped with p-type impurities generally formed from the side of a base substrate."}
{"text": "(1) Field of the Invention\nThe present invention relates to a method for processing a silver halide color photographic material and a color photographic developing composition used in the method, and more particularly to a method for processing a silver halide color photographic material in which the stability and the color-forming property of a color photographic developing solution are improved, and the increased fogging problem in continuous processing is lessened; and a color developing composition which can afford the color developing solution.\n(2) Description of the Prior Art\nColor developing solutions containing an aromatic primary amine color developing agent have long been used for forming color images, and they now play a main role in methods of forming color photographic images. However, it is well known that these color developing solutions have the disadvantage of oxidizing easily with air or metals, and if color images are formed using oxidized color developing solution there is increased fogging or the sensitivity or gradation changes, thus interfering with the desired photographic properties.\nAccordingly, hitherto various means of improving the preservative property of color developing solutions were studied, and, in particular, means using both hydroxyl amine and sulfite ions are most generally practiced. However, since hydroxyl amine when decomposed releases ammonia, which causes fogging, and since sulfite ions disadvantageously hamper, for example, the color-forming property when used as a competing compound in a developing agent, it is difficult to consider them preferable compounds (preservatives) to improve the preservative property of color developing solutions.\nAlthough sulfite ions, in particular, have long been used to improve the preservative property of various developing agents or to prevent hydroxyl amines from decomposing, they greatly hamper the color-forming property and markedly lower the color density. Especially when the sulfite ion is used in a system free of benzyl alcohol which is desirable to avoid the solution-preparation and environmental pollution problems involved with benzyl alcohol.\nAs compounds that can substitute for sulfites were suggested alkanolamines, described in Japanese Patent Application (OPI) No. 3532/1979, and polyethyleneimines, described in Japanese Patent Application (OPI) 94349/1981. Even if these compounds are used, however, they do not achieve enough desirable effects.\nTo improve the stability of color developing solutions, various preservatives and chelating agents have hitherto been studied. For example, as preservatives can be mentioned aromatic polyhydroxy compounds described, for example, in Japanese Patent Application (OPI) Nos. 49828/1977, 160142/1984, and 47038/1981 and U.S. Pat. No. 3,746,544, hydroxycarbonyl compounds described in U.S. Pat. No. 3,615,503 and British Patent No. 1,306,176, .alpha.-aminocarbonyl compounds described in Japanese Patent Application (OPI) Nos. 143020/1977 and 89425/1978, metal salts described in Japanese Patent Application (OPI) Nos. 44148/ 1982 and 53749/1982, and hydroxamic acids described in Japanese Patent Application (OPI) No. 27638/1977. As chelating agents can be mentioned amino polycarboxylic acids described in Japanese Patent Publication Nos. 30496/1973 and 30232/1969; organic phosphonic acids described in Japanese Patent Application (OPI) No. 97347/1981, Japanese Patent Publication No. 39359/1981, and West German Patent No. 2,227,639; phosphonocarboxylic acids described, for example, in Japanese Patent Application (OPI) Nos. 102726/1977, 42730/1978, 121127/1979, 126241/1980, and 65956/1980; compounds described, for example, in Japanese Patent Application (OPI) Nos. 195845/1983 and 203440/1983 and Japanese Patent Publcation No. 40900/1978; and organic phosphonic acid type chelating agents described in Research Disclosure Nos. 18837 and 17048.\nSatisfactory results, however, have not yet been obtained, because if organic phosphonic acid type chelating agents of the present invention and the above technique are used, the preservative property is not enough and the photographic characteristics are adversely affected.\nFurther, it is described in Japanese Patent Application (OPI) Nos. 95345/1973 and 232342/1984 that color photographic materials containing a silver chlorobromide emulsion with a high chlorine content are liable to cause fogging when the materials are color-developed. If such an emulsion is used it is essential to use a preservative that dissolved less in the emulsion and has better preservative performance, and in this sense satisfactory preservatives have not yet been found."}
{"text": "The invention relates to a side strut for a lower steering arm of a three-point attaching device of a tractor. The strut includes a first articulating telescopic member secured to the tractor and a second articulating telescopic member secured to the lower steering arm. The two telescopic members are adjustable relative to one another along an axis of adjustment between an inserted position and a maximum extracted position. The side strut further includes a locking pawl which is attached to one of the two telescopic members. The pawl is pivotable around an axis intersecting the axis of adjustment at right angles. The pawl includes two parallel cheeks arranged at the sides of the two telescopic members and connected to one another by a bridge. The side strut also includes two locking stops which extend transversely relative to the axis of adjustment. The stops project from the other of the two telescopic members and are form-fittingly embraced by locking recesses when the locking pawl is in a locking position.\nSide struts are known from DE 41 18 683 C1, issued May 27, 1992. In the case of the embodiment described therein, the movement between the telescopic members is limited by a longitudinal slot in the two cheeks of the locking pawl."}
{"text": "1. Field of the Invention\nThe present invention relates to a fully-automated washing machine for performing a laundry washing.\n2. Description of the Prior Art\nGenerally, in a conventional washing machine, as illustrated in FIG. 1, a body 1 of a washing machine is disposed with a tub 5 having a washing/spin-drying drum or tank 3 therein and a decelerater 9 is disposed underneath a center of the tub 5 for agitaing the laundry tossed into the washing/spin-drying drum by rotating clockwise and counter-clockwise the stirring blades 7 of a pulsator and the like installed on a floor within the drum 3.\nThe decelerator 9 is designed to receive a power in order to rotate the stirring blades 7 from a driving means 11 fixed to one end of a lower external side of the tub 5 through a power transfer means 13.\nA drainage means 16 such as a solenoid valve and the like is disposed on a middle of a drainage pipe 15 underneath the tub 5 in order to discharge the laundry water in the washing/spin-drying drum to the outside.\nFurthermore, on top surface area la of the body 1 of the washing machine, there is a lid 2 for covering an opening for insertion and removal of the laundry, and at the same time, on a rear side of the lid 2, in other words, on a lower rear side of a control panel 17, there is a detergent intake device 19 for dissolving and discharging a detergent into the drum 3.\nIn the washing machine thus constructed, when it is used by opening the lid 2 formed on the top surface area 1a of the body 1 and tossing the laundry into the drum 3, the water is supplied into the drum 3 and at the same time, the dissolved detergent is thrown into the drum from the detergent intake device 19.\nWhen a predetermined quantity of water and detergent is supplied into the drum 3, the driving means 11 starts to operate. At this moment, a turning effect of the driving means 11 is transferred to the decelerator 9 through the power transfer means 13 and pulleys 12 and 14 to thereby rotate the decelerator 9 in opposite directions, so that the washing is performed by way of the currents of water formed by the oscillation of the stirring blades 7 attached to a driving axis of the decelerator 9.\nMeanwhile, after the washing is finished, in accordance with an operation of the drainage means 16 attached to the drainage pipe 15, the water is discharged to the outside of the body 1, and in sequence, a washing cycle comprising another intermediate spin-dry, rinse, spin-dry and the like is automatically executed to thereby perform the laundry washing.\nHowever, in the washing machine thus constructed to perform the washing by way of the water currents, organic or inorganic materials attached to the laundry can be separated to some extent, but insoluble organic or inorganic materials such as dirt, fat and the like can hardly be separated from the laundry solely by water currents.\nTherefore, these kinds of insoluble foreign objects are separated by the detergent, in other words, by an assistance of a surface active agent for changing respective surface tensions of the water, fabrics and the foreign objects, to thereby be washed out. However, the foreign objects stuck to underwears, socks and other white-colored fabrics can not be washed thoroughly with the detergent, so that such laundry is first boiled in a separate container.\nHowever, because a boiled washing method thus described is performed after the laundry, the detergent and the water are poured into a separate container to which heat is applied by gas, electricity and the like, a special attention should be always paid to avoid a burning of the laundry or danger of fire."}
{"text": "1. Field of the Invention\nThe present invention relates to a vacuum processing system and a method of operating thereof, and more particularly to a vacuum processing system having an exhaust system configured to exhaust the vacuum processing system and optionally exhaust a gas injection system coupled to the vacuum processing system.\n2. Description of Related Art\nTypically, during materials processing, when fabricating composite material structures, a plasma is frequently employed to facilitate the addition and removal of material films. For example, in semiconductor processing, a dry plasma etch process is often utilized to remove or etch material along fine lines or within vias or contacts patterned on a silicon substrate. Alternatively, for example, a vapor deposition process is utilized to deposit material along fine lines or within vias or contacts on a silicon substrate. In the latter, vapor deposition processes include chemical vapor deposition (CVD), and plasma enhanced chemical vapor deposition (PECVD).\nIn PECVD, a plasma is utilized to alter or enhance the film deposition mechanism. For instance, plasma excitation generally allows film-forming reactions to proceed at temperatures that are significantly lower than those typically required to produce a similar film by thermally excited CVD. In addition, plasma excitation may activate film-forming chemical reactions that are not energetically or kinetically favored in thermal CVD. The chemical and physical properties of PECVD films may thus be varied over a relatively wide range by adjusting process parameters.\nMore recently, atomic layer deposition (ALD), and plasma enhanced ALD (PEALD) have emerged as candidates for ultra-thin gate film formation in front end-of-line (FEOL) operations, as well as ultra-thin barrier layer and seed layer formation for metallization in back end-of-line (BEOL) operations. In ALD, two or more process gases, such as a film precursor and a reduction gas, are introduced alternatingly and sequentially while the substrate is heated in order to form a material film one monolayer at a time. In PEALD, plasma is formed during the introduction of the reduction gas to form a reduction plasma. To date, ALD and PEALD processes have proven to provide improved uniformity in layer thickness and conformality to features on which the layer is deposited, albeit these processes are slower than their CVD and PECVD counterparts."}
{"text": "Electronic devices connected via a network typically require electrical connectors to be connected to the network or other relevant electronic devices.\nA conventional electrical connector typically requires an adapter terminal to realize electrical connection between a PCB board and an input terminal or an output terminal, so it has a complicated structure, and requires complicated assembling processes, which leads to a poor network transmission performance.\nThe above information disclosed in the BACKGROUND is only for a better understanding of the background of the present disclosure, and therefore, it may include information that does not constitute known prior art to those skilled in the art."}
{"text": "1. Field of Endeavor\nThe present application relates to gas spectroscopy and more particularly to a fiber optic coupled multipass gas minicell.\n2. State of Technology\nThis section provides background information related to the present disclosure which is not necessarily prior art.\nU.S. Pat. No. 8,309,929 for tunable photonic cavities for in-situ spectroscopic trace gas detection includes the following information: “Gas analysis is conventionally performed using laboratory analytical techniques, e.g., gas chromatography or mass spectrometry (GC-MS), which do not satisfy current device and material constraints for unattended, flexible ground sensors, or for lightweight, highly sensitive systems for avionic operations. Absorption spectroscopy is a powerful alternative approach for gas in-field detection and identification, and several interesting techniques have been developed including tunable diode laser absorption spectroscopy (TDLAS). Typically, this occurs in the infrared (IR) region of the spectrum. Recently, micromechanically tunable vertical-cavity surface-emitting lasers (VCSELs) have been implemented in such fashion for near infrared (NIR) spectroscopy. Unfortunately, many existing TDLAS systems exhibit drawbacks that limit their deployment, including the need for cryogenic cooling, a requirement for a bulky multipass cell, or a long hollow or porous fiber with a relatively slow time response.”\nTunable Absorption Spectroscopy is a common technique for gas detection and simply provides the chemical fingerprint of components by the absorption of the light at the wavelengths which resonate with the molecules inducing vibrations and rotations. By simply shining light through a volumetric gas sample on a detector, when the wavelengths is right, quench will occur and a dip in the power will be recorded at the detector.\nIn order to improve sensitivity various approaches are taken, from increasing the number of passes of light through the volume (thus increasing the absorption in discrete defined spaces) to implementing modulation techniques to remove background and noise.\nApplicants have developed an extremely small and robust cell in which a very small volumes of gas is sampled while maintaining high sensitivity and specificity by combining it with highly tunable VCSELs (to provide various absorption lines of one specie and capture varies elements at once) and wavelength modulation spectroscopy (to extract high harmonics which are not affected by the 1/f noise)."}
{"text": "(a) Field of the Invention\nThe present invention relates to a method of forming a trench for semiconductor device isolation and, more particularly, to a method of forming a shallow trench for semiconductor device isolation which exhibits good performance characteristics, and is well adapted for serving in device miniaturization.\n(b) Description of the Related Art\nGenerally, a local oxidation of silicon (LOCOS) technique has been used for device isolation purpose in the semiconductor fabrication process.\nIn the LOCOS technique, the silicon wafer itself is thermally oxidized while using a nitride for a mask, and requires only a small number of processing steps. For this reason, the resulting oxide involves lower degree of stress, and exhibits good insulating characteristics.\nHowever, in the application of the LOCOS technique, the device isolation area takes a large volume that limits device miniaturization as well as involves occurrence of the so-called bird\"\"s beak.\nIn order to solve the above problems, a shallow trench isolation (STI) technique has been suggested.\nIn the STI technique, a shallow trench is made in the silicon wafer, and filled with an insulating material. In this structure, the device isolation area is small, and hence the resulting trench can be well adapted for device miniaturization.\nFIGS. 1A to 1D schematically illustrate the steps of processing a trench for device isolation according to a prior art.\nAs shown in FIG. 1A, a silicon wafer 1 is thermally oxidized so that a pad oxide 2 with a thickness of 100-200 xc3x85 is grown at the silicon wafer 1. A nitride 3 with a thickness of 1000-2000 xc3x85 is deposited onto the pad oxide 2 through low pressure chemical vapor deposition (LPCVD). The pad oxide 2 and the nitride 3 are patterned through photolithography, and the exposed silicon wafer 1 is etched such that a shallow trench is formed at the device isolation area.\nAs shown in FIG. 1B, the silicon wafer 1 is thermally oxidized so that a liner oxide 4 is grown at the inner wall of the trench. A trench-filling oxide 5 with a thickness of 8000-11000 xc3x85 is deposited onto the entire surface of the silicon wafer 1 through atmospheric chemical vapor deposition (APCVD) such that the trench is completely buried by the trench-filling oxide 5. The silicon wafer 1 is then annealed at 900-1000 xc2x0 C. for 20-50 minutes under a nitrogen gas (N2) atmosphere to thereby density the trench-filling oxide 5.\nAs shown in FIG. 1C, the densified trench-filling oxide 5 is patterned through photolithography such that it is left only at the trench area.\nAs shown in FIG. 1D, the trench-filling oxide 5 at top side of the trench is removed through chemical mechanical polishing (CMP) while utilizing the nitride 3 as a stopping layer during the polishing. That is, the trench-filling oxide 5 is planarized such that the top portion of the trench-filling oxide 5 is level with the nitride 3. In this way, the formation of a shallow trench for device isolation is completed.\nHowever, the above-described trench formation technique bears complicated processing steps. Furthermore, when the trench-filling oxide 5 is planarized through the CMP, the top edge portion of the trench-filling oxide 5 cracks so that the trench-filling oxide 5 remaining after the planarization is torn (usually called the xe2x80x9ctorn oxidexe2x80x9d), causing serious device failures.\nFurthermore, among the impurities used in the CMOS transistor, boron-like components that exhibit a high diffusion property at high temperatures are diffused toward the trench during annealing or thermal oxidation and can induce an electrical short or leakage, deteriorating device reliability.\nIt is an object of the present invention to provide a method of forming a trench for device isolation which has simplified processing steps, avoids the tom oxide phenomenon, and ensures device reliability.\nThese and other objects may be achieved by a method of forming a trench for semiconductor device isolation including the following steps. A trench is first made at a device isolation area of a silicon wafer by etching the silicon wafer through a mask pattern. A liner oxide is then formed on the silicon wafer and within the trench through thermal oxidation, and a nitride layer is formed on the liner oxide through low pressure chemical vapor deposition. The nitride is anisotropically dry-etched such that the nitride is left only at the sidewall of the trench. A trench-filling oxide is then deposited onto the entire surface of the silicon wafer through high pressure chemical vapor deposition, and annealed. The trench-filling oxide is planarized through chemical mechanical polishing until the top surface of the trench-filling oxide layer is positioned slightly over the liner oxide on the silicon wafer. The silicon wafer is then wet-cleaned, and thermally oxidized such that a pad oxide is grown at the surface of the silicon wafer.\nThe thickness of the liner oxide is established to be 150-400 xc3x85, and that of the nitride to be 300-1000 xc3x85. The thickness of the pad oxide is 100-200 xc3x85."}
{"text": "1. Field of the Invention\nThe invention relates to a keyboard apparatus used in electronic keyboard instruments such as an electronic organ, an electronic piano, and a synthesizer, or in electric keyboard instruments.\n2. Description of the Related Art\nAs a keyboard apparatus used in electronic keyboard instruments such as electronic organs and electronic pianos, there has conventionally been provided a comb-teeth shaped keyboard apparatus such as that disclosed in, for example, JP H06-342281 A. This keyboard apparatus is structured such that a plurality of key main bodies to which a key depression operation is performed are coupled via thin coupling parts to a common key support part extending in a key arrangement direction and are supported to be pivotable in a key depression direction, and they are integrally formed of resin to form a comb-teeth shaped key unit. For example, in each one octave key range, a white key unit consisting of four alternate white keys, a white key unit consisting of the remaining three white keys, and a black key unit consisting of five black keys are separately formed, the respective common key support parts are integrally stacked, and the integrated key unit is fixed to a key mounting part of a keyboard frame with screws. Such integrated key units in number corresponding to a necessary key range are connected, whereby a keyboard apparatus is structured.\nWith this structure, since it is not necessary to provide fulcrum parts in the keyboard frame and the plural keys are integrated as a unit, the less number of components are required and the assembly maintenance are easy, which is extremely effective for reducing manufacturing cost.\nHowever, it has been necessary to provide a key guide for each key in order to restrict the key arrangement direction position of each key and to prevent yawing at the time of key depression.\nTherefore, there has been proposed a comb-teeth shaped keyboard apparatus requiring no key guide as disclosed in JP H07-92963 A. In the keyboard apparatus, a total width in a key width direction of a coupling part in each key unit as described above is larger than a width of a rear end of a key main body, thereby realizing sufficient flexibility in a key depression/release direction yet increasing a sectional secondary moment against a force in the key width direction (lateral direction). Accordingly, it is possible to restrict yawing in the key width direction of the key main bodies without using any key guide, resulting in reduced manufacturing cost of the keyboard apparatus.\nFurther, a keyboard apparatus of an electronic musical instrument as disclosed in, for example, U.S. Pat. No. 5,834,668 B has been conventionally in use. In this keyboard apparatus, in order to provide a heavy key touch feeling comparable to a key touch feeling of an acoustic musical instrument such as a piano, mass members such as hammers are provided for respective keys. Each of these hammers pivots in linkage with a depression operation of each key to produce a force depending on its movement, and the force works as a counter force against a key depression force, thereby providing a desired key depression touch feeling.\nHowever, the comb-teeth shaped keyboard apparatus as disclosed in JP H06-342281 A requires the key guides for the respective keys, which has posed a problem that the structure of the whole keyboard apparatus cannot be sufficiently simplified and it takes time to adjust the positions of the key units and the key guides.\nHere, structuring the key units as disclosed in JP H07-92963 A maybe disuse the key guides. However, it is necessary to make the coupling parts sufficiently wide in the key width direction so that the coupling part of each key partly overlap in plane with the coupling parts of the adjacent keys when the keyboard apparatus is assembled by stacking the common key support parts of the plural key units. Otherwise, key guides need to be provided on a key free end side. Therefore, the shape of the key units and the structure of the keyboard become complicated, resulting in increased molding cost and design restrictions.\nMoreover, in a case of a mini keyboard whose keys are shorter than keys in a standard keyboard, such a structure with the wide coupling parts can sufficiently restrict yawing of key main bodies, but in a case of a standard keyboard, it is difficult to sufficiently restrict yawing of key main bodies (especially, white keys) only by the wide coupling parts since the key main bodies are long, which poses a problem of bad performability such as difficulty in playing glissando.\nFurther, in the keyboard apparatus having the mass members operating in linkage with the respective keys as disclosed in U.S. Pat. No. 5,834,668 B, a large number of members such as key guides and key switches for the respective keys, mechanisms for linking the keys and the mass members (hammers), and upper/lower limit stoppers of the mass members have to be provided under the keys, which imposes many restrictions on the shape, arrangement, movable range, and so on of the mass members, and therefore it has been extremely difficult to design the keyboard apparatus providing a desired touch feeling and being low cost."}
{"text": "1. Field of Technology\nThe present application relates to a device for locking an equipment item onto a locking track attached to the floor of a cabin, including two parallel retaining edges divided by a longitudinal slot and a number of spatially separated openings, wherein the device has on its upper side means for locking the equipment item and on its lower side includes at least two locking projections the size and the arrangement of which are adapted to at least two adjacent openings, and wherein by shifting the locking projections in the longitudinal direction of the locking track until the locking projections are located between two adjacent openings, a positive connection between the device and the locking track can be produced.\n2. Brief Description of Related Art\nSuch a device is known from U.S. Pat. No. 5,871,318, for example, and is used in particular for the longitudinal shifting of aircraft seats along locking tracks that are attached onto the cabin floor of an aircraft. In order to fasten the aircraft seats onto the locking tracks, the device is initially aligned such that a support component with molded projections is aligned with the circular openings, so that this can be moved through the longitudinal slot in the locking track. The support component is then moved by half the distance between adjacent openings in the longitudinal direction of the track, wherein the T-shaped cross-sectional projections engage behind the retaining edges on both sides of the locking track and prevent the support component from coming out of the locking track. In a third step, locking parts having a corresponding contour are then moved down by the support component, so that these are retained in the openings which are now located below. In this context, the locking parts are pushed against the retaining edges, which results in locking the support component and thus the seat attached to it in the track.\nIn this case it is disadvantageous that the support components, after having been inserted into the locking track, must be moved as accurately as possible by half the distance of the opening, so that the locking parts engage into the openings and thus making secure locking possible. Since several seats are normally combined in one row and are arranged on one seat support, it is necessary that this entire row of seats must be moved. This process requires some time for assembly, particularly if the distance by which the row of seats is shifted deviates too greatly from the setpoint, and locking is not possible initially until some fine alignment has been done. This shifting process can also result in some tilting and can thus produce jamming."}
{"text": "Floor lamps, desk lamps and table lamps are all very well known in the lighting industry. Many of such lamps contain multi-position switches for changing the level of illumination provided by the general area lighting means incorporated within such lamps. Low light level lamps, commonly referred to as night lights, are also well known in the lighting industry and are constructed having many different configurations. It is also known in the prior art to incorporate a low light level light bulb into a table lamp, floor lamp or desk lamp. In many instances, the incorporation of the low light level light bulb accomplishes the additional function of providing a more decorative appearance to the body of the lamp. The present invention is specifically directed to a lamp, which can be a desk lamp, table lamp or floor lamp, of the type which also incorporates as an integral part thereof one or more low light level bulbs housed within the body of the lamp to provide decoration as well as a night light function. In such prior art structures, it is extremely difficult to replace the low light level light bulb when it has burned out. Prior art lamps of this type generally require extensive dismantling of the lamp to gain access to the interior of the lamp to replace the low light level light bulb. In one such prior art lamp the base had to be removed to replace the bulb. One lamp which solves the foregoing problems is shown in U.S. Pat. No. 6,916,108 by the same inventor and assigned to the same assignee as this application and the present invention is an improvement thereon.\nTherefore, there is need in the lighting industry for a lamp which incorporates both a general area lighting means and low light level light bulbs (night light) which also provides a means for ready access to the interior of the body so that the night lights may be easily changed when such is required."}
{"text": "This invention is concerned with the improvements in and relating to an information recording method of recording video information on an information storage medium and an information reproducing method of reproducing the video information from the information storage medium, and more particularly to those suitable for a case where the video information recorded on the information storage medium is the digital video information compressed according to the MPEG (Moving Picture Image Coding Experts Group) standards.\nIn recent years, systems for playing back an optical disk on which video (or moving picture) information and audio information have been recorded have been developed. They have been widely used in the form of, for example, LDs (Laser Disks) or video CDs (Compact Disks) for the purpose of reproducing movie software or karaoke.\nIn this connection, the DVD (Digital Versatile Disk) standard employing the internationally standardized MPEG-2 scheme and the AC (Audio Compression)-3 or other audio compression schemes has been proposed. The DVD standard covers playback-only DVD video (or DVD-ROM (Read-Only Memory)), write-once DVD-R (Recordable), rewritable DVD-RAM (Random Access Memory) (or DVD-RW (Rewritable)).\nThe DVD video (DVD-ROM) standard supports MPEG-2 for moving picture compression scheme and not only liner PCM (Pulse Code Modulation) but also AC-3 audio and MPEG audio for audio recording scheme.\nThe DVD video standard further supports sub-picture data obtained by run-length compressing the bit map data for subtitles and reproduce control data (navigation data) for data searching by fast-forward playback or fast-rewind playback.\nFurthermore, the DVD video standard supports ISO (International Organization for Standardization) 9660 and UDF (Universal Disk Format) to allow computers to read data.\nFor DVD video (DVD-ROM) optical disks, a 12-cm diameter single-sided single-layer disk has a storage capacity of about 4.7 GB (Giga Bytes); a 12-cm diameter single-sided double-layer disk has a storage capacity of about 9.5 GB; and a 12-cm diameter double-sided double-layer disk has a storage capacity of about 18 GB, provided that 650-nm (nanometers) wavelength laser light is used for reading.\nOn the other hand, for DVD-RAM (DVD-RW) optical disks, at the present time, a 12-cm diameter single-sided disk has a storage capacity of about 2.6 GB and a 12-cm diameter double-sided disk has a storage capacity of about 5.2 GB. Namely, DVD-RAM optical disks in practical use have a smaller storage capacity than DVD-ROM disks of the same size.\nIn playback-only DVD video (DVD-ROM), like a hierarchical file structure used by a general-purpose computer operating system, the directory structure of information (data files) recorded on an information storage medium is such that a subdirectory of video title set VTS and a subdirectory of audio title set ATS are connected to a root directory as shown in FIG. 1.\nIn the subdirectory of video title set VTS, various video files (including VMGI, VMGM, VTSI, VTSM, and VTS) are so arranged that the individual files can be managed in order. A specific file (for example, a specific VTS) can be accessed by specifying a path from the root directory to the file.\nSpecifically, the root directory of a DVD video disk includes a subdirectory called video title set VTS. The subdirectory can contain various management data files including VIDEO_TS.IFO or VTS—01—0.IFO, backup files, including VIDEO_TS.BUP and VTS—01—0.BUP, for backing up the information in those management data files, and a video data file VTS—01—1.VOB managed on the basis of the contents of the management data files and used to store digital video information. The subdirectory can also contain menu data files (including VMGM and VTSM) for storing specific menu information.\nA DVD video disk is composed of a video manager VMG and at least one or up to 99 video title sets VTSs. The video manager VMG is composed of control data VMGI, VMG menu video object set VMGM_VOBS, and backup control data VMGI_BUP. Each data is recorded on an information storage medium as a single file.\nAs shown in FIG. 1, on the DVD video disk, the individual video title sets (e.g., video title set VTS #1 and video title set VTS #2) have to be recorded in separate files. In each video title set (e.g., video title set VTS #1), control data VTSI, VTS menu video object set VTSM_VOBS, and backup control data VTSI_BUP are recorded in separate files. Additionally, title video data VTS—01—1.VOB and VTS—01—2.VOB in the VTS are recorded in plural files.\nThe DVD-RAM disk uses a UDF file system, not a FAT (File Allocation Table) file system. The details of UDF will be described in details later. Like FAT, UDF enables a hierarchical structure of files and records data in files on an information storage medium. In the prior art, both of the UDF file and the FAT file are filled with data and have no unrecorded area in them.\nThe contents will be explained in detail using one example. For example, when a statement has been written using word processor software (such as Ichitaro, Word, or Amipro) running on a PC (Personal Computer), the written statement is recorded on an information storage medium as a file. In this case, all the file is filled with text data. Even if a space area or a continues enter mark portion with no sentence continues long in the middle of the written sentence, that portion in the stored file will be filled with space data and enter data and therefore there will be no fully unrecorded area in the file.\nEven when the user reads the document file and stores the data after deleting the middle of the sentence, an unrecorded area is never defined in the stored information and is recorded on the information storage medium as a file with the data items before and after the deleted portion putting together. As a result, the size of the file recorded on the information storage medium decreases by the amount of data in the deleted portion.\nWith application software running on an ordinary PC, a file read from an information storage medium for editing is transferred as it is to a buffer memory (semiconductor memory) on the PC. The edited data is stored temporarily in the buffer memory on the PC. Once the user has given an instruction to store the file, the edited data stored in the buffer memory on the PC is written over the whole file on the information storage medium. As described above, with the conventional file system, such as a FAT or UDF file system, when the file data is changed, all the data in the file is changed at a time in the overwrite process. This is different from the present invention where the data in only a part of the file is changed.\nFIGS. 2A and 2B illustrate examples of reproducing video information using program chains PGCs on a DVD video disk. As shown in FIG. 2A, the playback data is divided into cells and playback sections from cell A to cell F are specified. In the individual program chains PGC #1 to #3, PGC information is defined as shown in FIG. 2B. Specifically, the table in FIG. 2B reads as follows.\n1. Program chain PGC #1 shows an example of being made up of cells specifying consecutive playback sections. The playback sequence is:\nCell A→Cell B→Cell C.\n2. Program chain PGC #2 shows an example of being made up of cells specifying intermittent playback sections. The playback sequence is:\nCell D→Cell E→Cell F.\n3. Program chain PGC #3 shows an example of being made up of cells specifying disorderly playback sections, regardless of the direction of playback or repetitive playback. The playback sequence is:\nCell E→Cell A→Cell D→Cell B→Cell E\nBy defining different program chains PGCs as described above, different display sequences can be realized for the same cells. In a DVD video disk, all the cell information is not necessarily displayed by a single program chain PGC because of the freedom of program chain PGC setting.\nWhat has been explained above is about the data structure of the video information recorded on a playback-only DVD video disk. An information storage medium capable of recording and reproducing video information using a DVD-RAM disk or a DVD-RW disk is now being developed as one form of the DVD family.\nIt is desirable that the video information recording format on the information storage medium capable of video recording and reproducing should have a continuity and a relation with the data structure of a DVD video disk. In addition, a UDF file system is used for DVD-RAM disks or DVD-RW disks, as in the playback-only DVD video disk.\nWhen the data structure of the aforementioned DVD video disk is used directly as the data structure on a recordable (videorecordable) information storage medium and the above conventional UDF (or FAT) file system is used, the following problems arise:\n1. Since the control data and video data are recorded in such a manner that they are distributed over plural files, when having deleted a file by mistake, the user is unaware of the position of the error until trying to reproduce the deleted file in the course of playback. For the playback-only DVD video disk, there is no possibility that the user will delete a file. In the case of recordable/erasable information storage mediums, however, there is a danger that the user will delete a file by mistake.\n2. Since the control data and video data are recorded in such a manner that they are distributed over plural files and the data structure has the same hierarchical structure as that of the computer data, it is difficult for the family user unfamiliar with the computer to understand the deleted place or the recorded place. Specifically, knowing only the VTR (Video Tape Recorder) as a medium capable of recording video information, the family user wonders which part of the single tape the place where the picture has been recorded or deleted has occupied. Therefore, showing the user small files of record or the result of erasing as they are would throw the user into confusion.\nAs shown in FIG. 1, in the DVD video disk, information is recorded in such a manner that it is divided into separate files by video title set VTS. Thus, when plural video title sets (VTS#1 and VTS#2 in FIG. 1) have been recorded on the information storage medium, the user familiar with only the VTR has no idea about the playback procedure.\n3. With a method of allowing the family user to select a specific cell corresponding to a program chain PGC for the recorded information, some user is liable to fall into confusion. Specifically, knowing only the VTR as a medium capable of recording video information, the family user will probably wonder which part of the single tape the place where the picture has been recorded or deleted has occupied. Therefore, it would be difficult for the user to understand the concept of selecting cells by a program chain PGC on a playback-only DVD video disk.\n4. In a data file recorded using the conventional UDF or FAT, there is no unrecorded area. Therefore, when part of a specific data item in a file has been deleted or a few pieces of video information have been added, the data items in front of and behind the deleted portion are squeezed together and connected or the pieces of information are added to the end of the existing data. Each time data is deleted or added, the size of the whole data file has to be changed and all the changed data file has to be recorded on the information storage medium again. As a result, it take a very long time to complete the editing process.\nSpecifically, with the conventional UDF or FAT, since a file has no unrecorded area, the following processes cannot be carried out:\n(a) The process of changing the erased place to an unrecorded area when part of the data in a file has been deleted.\n(b) The process of recording additional data on an unrecorded area in a file without changing the entire file size.\nTherefore, each time the data is deleted partially or added, the file size has to be changed.\nAs a result, the entire file has to be recorded again on the information storage medium. In the case of a video file in which video information has been recorded, the size of a single video file is as large as more than several hundreds of megabytes (MB). If a file as large as several hundreds of megabytes is all recorded again on the information storage medium each time a slight change has been made, it will take an extremely long time to change the contents of the file.\nFor a further description of the prior art, reference may be made to:\nJapanese Patent Application No. 040876 (filed on Feb. 23, 1998)\nJapanese Patent Application No. 040877 (filed on Feb. 23, 1998)\nJapanese Patent Application No. 040879 (filed on Feb. 23, 1998)."}
{"text": "Conference calls are commonly used to connect a group of participants together through voice and/or video data in real time. The quality of the conference call data is important for capturing accurate data to reflect the discussions performed during the conference call. For example, a conference call may be recorded and the voice recording may be used to generate a textual transcription of the words spoken by the participants during the conference call.\nConventionally, a transcriptionist may be present during the call to hear and type the words spoken by the participants. Other ways to capture the content of the conference call may be to record the voice and use a digital software program application to convert the voice or speech into text. Voice recognition software programs may be used to convert human speech into text so that the words spoken by the conference call participants may be later referenced by others to review the content of the conference call.\nCurrently, a conference call is recorded by a single recording effort, which is stored in a single recorded file. A transcription of the recording may then be enacted via a speech-to-text software application that identifies the words spoken by the conference call participants and converts the words into text. Another option may be to use a human transcriptionist, or a combination of both a human transcriptionist and a software application to convert the speech into text.\nCurrent software applications are limited to a single recording of a conference call with limited ability to distinguish between different speakers. Such applications have a low percentage of accuracy due to the limited capabilities of the software application. For example, the software application may not be capable of distinguishing the different voices from one another. Also, the different recorded voices may be overlapping, and some voices may be louder than other recorded voices. These various shortcomings of a single data recording may cause interference when attempting to transcribe the audio content of the voices of the conference call participants."}
{"text": "In the current market of products for infusion, such as coffee, barley coffee, tea, camomile and the like, the use of single-dose “pods” has increased considerably and nowadays a very popular way of making hot beverages is to use such pods in specially designed machines, even for household or office use (that is, for small to medium quantities).\nNormally, these pods are of the single-dose type, used to brew a single serving of the beverage. This type of pod, to which the present specification refers but without thereby restricting the scope of the inventive concept, usually consists of two portions of filter paper placed one over the other and sealed to enclose a single product dose of substantially circular shape.\nA prior art packaging machine for making pods of this kind is described and illustrated in European Patent No. EP 432.126, where the pods are formed according to a process which comprises the steps of: feeding a first web or length of filter paper to a station where suitable means cause the filter paper to be wrinkled or crinkled; moving the web of filter paper along the surface of a forming drum, provided with circular pockets and with suction means which force spaced areas of the filter paper into the pockets in the drum in such a way as to form a succession of substantially circular pouches in the filter paper; filling a dose of product into each pouch by means of a dosing station located downstream of the suction drawing belt in the direction of rotation of the pouch forming drum and consisting of a second revolving drum synchronised with the pouch forming drum; and, lastly, joining the first web of filter paper, with the product-filled pouches in it, to a second web of filter paper fed at a respective sealing station located downstream of the filling station, again relative to the direction of rotation of the forming drum.\nThe method embodied by the machine described in the aforementioned European Patent has several major disadvantages, especially in connection with the steps of drawing the filter paper by suction to form the pouches.\nIndeed, although the method includes a complex step of preparing the web of filter paper in a specific station that crinkles the filter paper to make it suitable for forming the pouches, many of the suction drawn pouches formed in the filter paper in this way are in actual fact faulty, in that they have creases or folds made in them. These creases not only prevent the infusion product subsequently filled into the pouch from spreading uniformly in the pouch itself but may also result in poor seals between the first and second webs.\nAs a result, the pods made are often defective and have to be rejected since they contain insufficient quantities of infusion product or have improperly sealed parts that create gaps through which the infusion product can escape.\nThe aim of the present invention is therefore to provide an apparatus for making pods for products for infusion that is free of the disadvantages of the prior art described above and that is capable of consistently making high quality pods."}
{"text": "Multicast communications have grown increasingly important and complex in networking architectures. Multicasting generally refers to a scenario in which multicast packets are routed to multiple destinations and generally replicated an appropriate number of times. Each particle of a multicast packet generally includes a reference count, which offers some information about the number of intended targets for the data. As a particle is created an element may operate to set the reference count, whereby once the particle's reference count reaches ‘0’ it is generally freed by the system. When a target identifies itself as ready to receive a portion of data, the data needs to be accessed and read from a memory location. In some environments, a problem is presented when a multicast packet or a data segment is stored in some memory unit and then recalled by multiple targets simultaneously. This may also present a problem in cases where the target does not yet know the correct value of the reference count or what other potential targets may also be attempting to access the data. The time interval required to adequately modify the reference count must be kept to a minimum, as such an operation generally inhibits system speed.\nIn cases where the target opts to wait to receive the data before modifying the reference count, another target may also access the data concurrently and thus retrieve an invalid reference count. In addition, other problems may be presented where the modified reference count is corrupted when written back to the memory unit by either the first target seeking access to the data or a subsequent target. These problems may result in corrupted data segments, inferior data communication speeds, inefficient uses of networking resources, and poor performance characteristics for associated multicasting operations."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method for conveying products and, more particularly, to an apparatus and method for conveying documents in a high-speed electronic finishing system.\n2. Description of the Prior Art\nIn the electronic finishing field, multi-page documents such as insurance policies and contracts are printed on high-speed laser printers and then finished on an electronic finishing system which assembles the document and adds any desired inserts (such as prefolded advertising brochures or other solicitations) to the document to create a product which are sometimes folded in manner set forth in U.S. Pat. No. 5,554,094. Then, the product is inserted into a standard envelope to create a package ready for mailing. These steps are done at high speed while the accuracy of the document assembly is verified through use of a printed bar code on each sheet of the document.\nThe prior art conveyors used to convey documents under the insert feeders are complicated and are difficult to adjust to accommodate different size documents. The prior art adjustable conveyors are inadequate because they often leave gaping holes in the conveyor path allowing products to sag or drop below the conveyor path.\nAnother problem with prior art conveyors is that the guiding rails are adjustable but the pusher path is fixed whereby, as the rails are adjusted wider for wider products, the pusher integrity become less stable. Thus, as the product gets wider in relation to the pusher path, the conveyor is the less reliable at advancing the product in a controlled manner.\nYet another shortcoming of the prior art conveyors is that the hold down devices are fixedly positioned such that they will prove adequate for the range of applications intended. However, when additional hold-down requirements arise, it often left up to the ingenuity of the operator or technician to solve the requirements.\nFinally, all of the components which overhang the prior art conveyors make them difficult to access to clear jams and clean.\nThe present invention is designed to overcome the above limitations that are attendant upon the use of xe2x80x9cprior artxe2x80x9d devices, and toward this end, it contemplates the provision of a novel apparatus and method for the high-speed conveying of products and the provision of inserts to be combined with the products as they are conveyed.\nIt is an object of the present invention to provide an apparatus and method utilizing adjustable spring loaded collapsible pusher fingers to move the product though the conveying system.\nYet another object is to provide an apparatus and method which accomplish the desired provision of inserts at high speeds.\nA further object is to provide an apparatus that is compatible with conventional electronic finishing systems and is generally compatible in physical size, form and configuration with such systems, to be readily adaptable for the same use without disadvantage.\nIt is a general aim of the invention to provide such an apparatus which may be readily and economically fabricated and will have long life in operation and significantly greater flexibility in use.\nIt has now been found that the foregoing and related objects can be readily attained in an apparatus and method in accordance with the present invention which conveys a product along a product conveying path, permits products of different widths to be conveyed and easily adds inserts as desired.\nAccording to t he invention, a conveyor system for conveying the product along the product conveying path utilizes six primary assemblies: pusher assembly, inner rail assembly, outer rail assembly, strap-ramp assembly, hold-down strap assembly and stop finger assembly. Insert feeders are spaced above the conveyor system to feed inserts onto the conveyor system so as to be combined with the product as the product is conveyed a long the product conveyor path.\nIn the pusher assembly, there are a plurality of pusher devices. Each pusher device has a plurality of pushers for engaging at railing end of the product to propel the product along the product conveying path. At least two of the pushers on each pusher device are spaced from one another transversely across the conveying path and moveable between first and second positions transversely across the conveying path. These pushers are closer to one another in the first position than in the second position thereby being adjustable to convey products of different widths. A mechanism is provided for moving the pusher devices so that the pushers move along the product conveying path to propel the product engaged at the trailing end thereof along the product conveying path.\nFor the inner and outer rail assemblies, pairs of inner opposed rails support the product as it is propelled along the product conveying path. The inner opposed rails are adjustable between at least two distinct positions to permit conveying products of different widths. Pairs of outer opposed rails cooperate with the pairs of inner opposed rails for supporting the product as it is propelled along the product conveying path. The outer opposed rails being adjustable relative to and independent of the pair of inner opposed rails to permit conveying of products of different widths. Each of the outer opposed rails having a base portion for supporting the product as it is propelled along the product conveying path and an upstanding portion for guiding the product as it is propelled along the product conveying path.\nWith the strap-ramp assembly, a strap support device has a plurality of straps to support the product being propelled along the product conveying path. The straps are spaced from one another transversely across the product conveying path. At least some of the straps are moveable between first and second positions. In the first position, these straps are located to support the product being propelled along the product conveying path. In the second position, these straps are not located so as to support the product being propelled along the product conveying path. Desirably, a mechanism is provided for moving these straps of the strap support device between the first and second positions.\nIn the hold down assembly, at least one hold down support device extends over the product conveying path and at least one hold down strap is removeably mounted on the hold down support device. The hold down device holds the product flat on the conveyor path and helps decelerate the product at each indexing movement of the product.\nTo stop the product accurately at each station following each indexing movement of the product, the stop finger assembly includes a stop finger support device cantilevered over the product conveying path with at least one stop finger removeably mounted on the stop finger support device. The stop finger support device is moved so that the at least one stop finger moves into the product conveying path to engaged the product at a leading end thereof following each indexing movement of the product and moves out of the product conveying path to begin each indexing movement of the product.\nTo feed at least one insert onto the conveyor system so as to be combined with the product as the product is conveyed along the product conveyor path, a plurality of insert feeders are spaced above the conveyor system. Conveniently, each of the insert feeders is pivotally mounted adjacent the product conveyor path whereby the insert feeder can be rotated between an operative position spaced above the conveyor system for feeding at least one insert onto the conveyor system so as to be combined with the product as the product is conveyed along the product conveyor path and a raised inoperative position spaced farther from the product conveyor path to provide access to the product conveyor path.\nThe invention will be fully understood when reference is made to the following detailed description taken in conjunction with the accompanying drawings."}
{"text": "Subscribers to broadcast television sometimes desire to view television programming at times other than when the programming is broadcast. While digital video recorders (DVRs) such as the TIVO® system are used to record and time shift television programs, they are cumbersome to set up and use, particularly where the subscriber has to interconnect the DVR with a television, set top box or other television accessories. In addition, while DVRs provide a benefit to consumers, they may reduce the effectiveness of broadcast commercials if consumers fast forward through the commercials when they are viewing recorded programming. A particular challenge for subscribers to Internet protocol television (IPTV) is that not all channels are simultaneously received by the subscriber as they are in over-the-air or cable television broadcast systems. Therefore, the subscriber's ability to record live television broadcasts is limited to only those channels they can receive at one time."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the packaging of electronic components. More particularly, the present invention relates to a shielded electronic component package and method of fabricating the same.\n2. Description of the Related Art\nAs the art moved to smaller, lighter weight, and higher frequency electronic devices such as cellular telephones, integrated circuit packages utilized in these electronic devices were increasingly placed closer to other electronic components and structures. Due to this reduced spacing, radiation such as electromagnetic or radio frequency (RF) radiation emanating from an integrated circuit package had a greater probability of interfering with the normal operation of an adjacent electronic component and vice versa. However, such interference was unacceptable depending upon the particular application.\nFurther, to prevent interfering with the normal operation of adjacent electronic devices, e.g., radio receivers, it was important to prevent the integrated circuit package from emanating unwanted radiation to the ambient air. Similarly, to prevent interfering with the normal operation of the integrated circuit package, it was important to protect the electronic component(s) of the integrated circuit package from radiation emanating from adjacent electronic devices, e.g., radio transmitters. Stated generally, it was important to prevent unwanted radiation from electronic component(s) of the integrated circuit package from reaching the surrounding air and vice versa.\nTo prevent unacceptable electromagnetic interference, a shielding system was used. Typically, an electrically conductive metallic enclosure was placed around the integrated circuit package after the integrated circuit package was mounted to the larger substrate such as the printed circuit mother board. However, fabricating such a metallic enclosure and separately attaching the metallic enclosure to the printed circuit mother board was relatively cumbersome, complex, and costly."}
{"text": "This disclosure relates to user interfaces (UI's) for touch input devices, and more specifically for touch input displays used with electronic devices such as tablets, smart phones, touch laptop computers, touch enabled display monitors and the like. Many devices provide user control over various operating parameters of the device. In many instances, a UI is operated by a user to increase or decrease a value or quantity associated with a software application the device is running. A UI may be provided for controlling volume or manipulating tone controls of an audio software application such as a streaming audio application. UI's are also used to control various other quantities where progressive increasing or decreasing of a quantity or variable, or directional scrolling through a list of data elements is desired. Previous examples of touch controls have difficulty in proving both gross (i.e. large) and fine grain control in a single control. All examples and features mentioned below can be combined in any technically possible way."}
{"text": "As a magnetic memory device using a magnetoresistive effect element of the above-mentioned kind, there has been known a magnetic random access memory (hereinafter also referred to as “the MRAM”). The MRAM stores information using a combination (parallel or antiparallel) of the respective magnetization directions of two ferromagnetic materials included in the magnetoresistive effect element. On the other hand, reading of the stored information is carried out by detecting a change in the resistance value of the magnetoresistive effect element (i.e. a change in electric current or voltage), which changes depending on whether the magnetization directions of the two ferromagnetic materials are parallel or antiparallel to each other.\nMRAMs currently operative make use of a giant magnetoresistive (GMR) effect. An MRAM using GMR elements capable of providing the GMR effect is disclosed in U.S. Pat. No. 5,343,422. In this case, the GMR effect is intended to mean a phenomenon that the resistance value becomes minimum when the magnetization directions of two magnetic layers parallel to each other along the easy axis of magnetization are parallel to each other, whereas when the magnetization directions of the two magnetic layers are antiparallel to each other, the resistance value becomes maximum. As MRAMs using the GMR element, there are a coercive force difference type (pseudo spin valve type) MRAM and a switching bias type (spin valve type) MRAM. The coercive force difference type MRAM uses a GMR element having two ferromagnetic layers and a non-magnetic layer sandwiched therebetween, and makes use of the coercive force difference between the two ferromagnetic materials to thereby have information written therein and read out therefrom. Now, assuming that a GMR element is composed e.g. of “nickel-iron alloy (NiFe)/copper (Cu)/cobalt (Co)”, the rate of change in resistance of the element assumes a small value of approximately 6 to 8%. On the other hand, the switching bias type MRAM uses a GMR element having a fixed layer whose magnetization direction is fixed by exchange-coupling with an antiferromagnetic layer, a sensitive magnetic layer whose magnetization direction is changed by an external magnetic field, and a non-magnetic layer sandwiched therebetween, and makes use of the difference between the respective magnetization directions of the fixed layer and the sensitive magnetic layer to thereby have information written therein and read out therefrom. For example, assuming that a GMR element is composed e.g. of “platinum manganese (PtMn)/cobalt iron (CoFe)/copper (Cu)/CoFe”, the rate of change in resistance of the GMR element assumes a value of approximately 10%. This value is larger than the rate of change in resistance of the coercive force difference type MRAM, but insufficient to further enhance the recording speed and the access speed of the GMR element.\nTo solve the above problems, there has been proposed an MRAM which uses magnetoresistive effect elements (hereinafter also referred to as “the TMR elements”) 120 constructed as shown in FIG. 11 for utilizing a tunnel magnetoresistive effect (hereinafter also referred to as “the TMR effect”), as storage cells. As shown in FIG. 12, this MRAM is comprised of a plurality of bit lines 105 arranged in parallel to each other, a plurality of write word lines 106 arranged in parallel to each other, and at the same time orthogonal to the bit lines 105, a plurality of read word lines 112 arranged along the write word lines 106, and a plurality of TMR elements 120 arranged at orthogonally crossing portions (intersections) where the bit lines 105 and the write word lines 106 orthogonally intersect with each other, in a manner sandwiched therebetween. In this case, as shown in FIG. 11, each TMR element 120 is comprised of a first magnetic layer 102, a tunnel barrier layer 103, and a sensitive magnetic layer 104 as a second magnetic layer, which are sequentially deposited in the mentioned order.\nIn this case, the TMR effect is intended to mean an effect that a tunnel current flowing through the tunnel barrier layer 103 varies with a relative angle between the respective magnetization directions of the first magnetic layer 102 and the sensitive magnetic layer 104 as two ferromagnetic layers sandwiching the tunnel barrier layer 103 as a very thin insulating layer (non-magnetic conductive layer). More specifically, the resistance value becomes minimum when the magnetization directions of the first magnetic layer 102 and the sensitive magnetic layer 104 are parallel to each other, whereas when they are antiparallel to each other, the resistance value becomes maximum. Further, in the MRAM making use of the TMR effect, when the TMR element 120 is composed e.g. of “CoFe/aluminum oxide/CoFe”, the rate of change in resistance of the MRAM assumes a high value of approximately 40%, and the resistance value thereof is large, and hence it is easy to match the MRAM with a semiconductor device, such as a MOSFET or the like, when they are combined with each other. This makes it possible to easily obtain a higher output compared with the MRAM including the GMR element, and enhancement of the storage capacity and the access speed is expected. The MRAM making use of the TMR effect stores information by changing the magnetization direction of the sensitive magnetic layer 104 of the TMR element 120 to a predetermined direction, using current magnetic fields generated by electric current caused to flow through the bit lines 105 and the write word lines 106. On the other hand, to read out information stored in the MRAM, electric current in a direction perpendicular to the tunnel barrier layer 103 is caused to flow through the TMR element 120 via the bit lines 105 and the read word lines 112, to thereby detect a change in the resistance of the TMR element 120. It should be noted that the MRAM utilizing the TMR effect is disclosed in U.S. Pat. No. 5,629,922, or in Japanese Laid-Open Patent Publication (Kokai) No. H09-91949."}
{"text": "The present invention relates to a clamp apparatus that clamps a workpiece or the like that is to be supplied for machining processing, and in particular relates to such a clamp apparatus in which it is arranged to perform clamp driving while multiplying the clamping force that acts on a piston member by fluid pressure.\nIn a clamp apparatus that employs fluid pressure, it is possible to strengthen the clamping force by increasing the pressure reception area of the piston member. However, in this case, the clamp apparatus is increased in size, and this is disadvantageous because it reduces the freedom in designing the system. Thus, various types of clamp apparatus have been implemented in which the clamping force is strengthened by the provision of a force multiplication mechanism (servomechanism). In particular, there are a number of examples of provision of a force multiplication mechanism when pressurized air is used as the pressurized fluid.\nIn the clamp apparatus described in FIG. 4 of JP Patent 5,129,378: a cylinder bore is formed within the clamp main body; a first piston fitted in the cylinder bore, and an axial second piston fitted into the first piston, are installed so as to slide freely; an output rod is formed integrally with the second piston; a tubular portion that projects from above into the interior of the upper end portion of the first piston is formed integrally with the clamp main body; and there is provided a force multiplication mechanism which multiplies force generated by fluid pressure acting on the first piston and transfers the result to the second piston. There is also provided a compression spring for maintaining clamping force to bias the first piston in the direction for clamping operation.\nThe force multiplication mechanism includes: a plurality of support holes in the tubular portion; a plurality of spheres installed in these support holes so as to be movable in radial directions; a part conical surface that is formed on the inner surface of the upper end portion of the first piston; and a plurality of cam grooves that are formed on the second piston so as to be capable of contacting against the plurality of spheres. A first fluid pressure operation chamber is defined below the first and second pistons in the cylinder bore, and a second fluid pressure operation chamber is defined above the first and second pistons."}
{"text": "This invention relates in general to reducing heat energy in semiconductor devices, and more particularly, to thermal conductors embedded within semiconductor chip structures for removing heat from one or more conductive circuit members thereof, wherein the conductive circuit members are disposed within or on a low k dielectric material having a low thermal conductivity. More specifically, the invention addresses the problem of removing heat generated by copper wiring within low dielectric glasses and nanofoams in semiconductor chip structures.\nModern Very Large Scale Integrated (VLSI) chips require cooling to improve reliability of their circuitry and interconnects, to increase their circuit switching performance, and to regulate thermally generated noise in their circuits. Cooling reduces the likelihood that a metal wire will form voids or a contact will become open. It also reduces the extent of time-dependent transistor mobility and threshold degradation which adversely affects circuit performance and operation. Furthermore, in typical complementary metal oxide semiconductor (CMOS) microprocessors, every reduction in temperature of 10xc2x0 C. produces a 2% rise in operating frequency. For CMOS transistors, high temperatures yield significantly larger leakage currents, due to the thermal generation of carriers. This deleterious current doubles every 11xc2x0 C. and is known to adversely affect the functional operation of dynamic and analog circuits.\nAs VLSI circuits shrink to improve performance and increase operating frequencies, higher amounts of heat are generated, for example, due to constant switching of these devices. The removal of heat within a semiconductor chip structure thus becomes a major obstacle to the efficient performance of the device. Therefore, a need continues to exist for enhanced heat removal techniques for semiconductor devices.\nIn addition to the continued reduction in chip size, new materials are being incorporated to increase circuit performance. For example, dielectric materials with a dielectric constant (k) lower than that of conventional oxide reduce the parasitic capacitance between neighboring conductors, thereby improving circuit speed. However, most low k dielectrics are also lower in density than silicon oxide, and exhibit lower thermal conductivities as well thus further increasing the need for enhanced heat removal techniques.\nForming conductive lines and vias for semiconductor connections can be accomplished by a variety of methods. Prior techniques have used deposition and subtractive etch of metal to form conductive lines, followed by deposition of a dielectric layer and etch of vias therein to connect conductive lines above and below the dielectric layer. Newer techniques include damascene and dual damascene processes. In these techniques, a dielectric is deposited and patterned with trenches for conductive lines (in a damascene process) and for lines and vias (in a dual damascene process). Metal is then deposited and any metal overlying the dielectric outside of the vias and the trenches is removed by a chemical mechanical polish (CMP).\nThe metal that has traditionally been employed in such processes is aluminum. However, today\"\"s chips are designed to run very fast and two effects limiting speed are conductive line resistance and RC coupling induced delay due to higher wiring density and closer spacing of conductive wires. Copper, which has a lower resistance, is replacing aluminum for wiring of semiconductor chip structures. RC coupling is being addressed by the use of lower dielectric constant (low k) dielectrics, such as porous silicas or polyimide nanofoams to replace conventional silicon dioxide dielectrics. However, the combination of copper and low k porous material creates a problem with heat dissipation. Because of the improved resistance and lower RC coupling that a copper/low k dielectric combination affords, more power per unit time is applied to the line, significantly increasing the heat dissipation requirements. These porous low k materials characteristically do not provide heat dissipation as well as high density glass dielectrics. Therefore, as the temperature of copper lines continues to increase due to current flow, resistance also increases, degrading device performance. To add to the problem, certain low k dielectrics, especially organic foam dielectrics, will degrade both structurally and electrically at temperatures around 350xc2x0 C. This temperature has been observed to be exceeded in copper wires on some products. Thus, there is a need for techniques which may be integrated with copper/low k dielectric semiconductor interconnection structures to improve the heat dissipation capability thereof so that device performance may be maintained optimal and consistent.\nTo summarize, therefore, the present invention comprises in one aspect a semiconductor chip structure which includes a substrate having first and second opposing surfaces, and at least one electrically conductive circuit member disposed above the first surface of the substrate. The at least one electrically conductive circuit member resides at least partially on a dielectric material, with the dielectric material being disposed between the at least one electrically conductive member and the substrate. The dielectric material has a low dielectric constant and a first thermal conductivity. The semiconductor chip structure further includes at least one thermal conductor thermally coupled to the at least one electrically conductive circuit member to provide a path for heat escape from the electrically conductive circuit member during operation. The at least one thermal conductor has a second thermal conductivity, and the second thermal conductivity is greater than the first thermal conductivity. A thermal sink is also provided thermally coupled to the second surface of the substrate, wherein the at least one thermal conductor facilitates heat transfer from the at least one electrically conductive circuit member to the thermal sink.\nIn a further aspect, a method of fabricating a semiconductor chip structure is provided. The method includes: providing a substrate having first and second opposing surfaces; forming at least one electrically conductive circuit member above the substrate on a dielectric layer, wherein the dielectric layer resides between the at least one electrically conductive circuit member and the substrate and has a low dielectric constant and a first thermal conductivity; disposing at least one thermal conductor thermally coupled to the at least one electrically conductive circuit member to provide a path for heat escape from the electrically conductive circuit member during operation, wherein the at least one thermal conductor has a second thermal conductivity, the second thermal conductivity being greater than the first thermal conductivity; and forming a thermal sink thermally coupled to the second surface of the substrate, wherein the at least one thermal conductor facilitates heat transfer from the at least one electrically conductive circuit member to the thermal sink.\nTo restate, semiconductor chip structures are described herein having one or more integrated thermal conductors for dissipating heat generated by conductive circuit members, for example wiring levels, disposed above the substrates of the chip structures. Each thermal conductor can comprise a thermally conductive plug or a stack of plugs which facilitate thermal coupling of heat produced by one or more conductive levels to the substrate of the chip structure and/or an upper surface of the structure. In one embodiment, the thermal conductors are placed in strategic locations within the semiconductor chip structures and either directly physically contact the circuit members to be cooled or pass close thereto so as to be thermally coupled to the members. As one example, the conductive circuit members reside on or within a first dielectric material having a first thermal conductivity and the stack of plugs comprise a second dielectric material having a second thermal conductivity, wherein the second thermal conductivity is at least 3xc3x97 greater than the first thermal conductivity. More particularly, the present invention solves the problem of cooling semiconductor chip wiring when dielectrics having low thermal conductivity are used for interlevel isolation. The invention can be employed within a structure using copper wiring in combination with low k dielectric glasses and nanofoams.\nAs an enhancement, a thermal sink is thermally coupled to a second surface of the substrate opposite the first surface above which the conductive circuit members are disposed. The at least one thermal conductor is positioned to facilitate heat transfer from the at least one electrically conductive circuit member to the thermal sink disposed over the second surface of the substrate. Preferably, one or more thermally conductive via structures are formed in the substrate integral with the thermal sink for further facilitating thermal transfer from the at least one thermal conductor into the thermal sink.\nAdditional features and advantages are realized through the techniques of the present invention. Other embodiments and aspects of the invention are described in detail herein and are considered a part of the claimed invention."}
{"text": "This invention relates generally to steam turbines, and more particularly to improved apparatus for controlling a flow of steam to such turbines.\nIn a steam turbine generator system, the turbine is normally maintained at a constant speed and steam flow is varied to adjust the torque required to meet the electrical load imposed on the generator. This type of control is provided by a main control system which varies the flow of steam to the high-pressure turbine, and in some instances to the low-pressure turbine, to meet the load demand. The main control system is designed to accommodate for normal changes in load demand and to smoothly adjust the turbine operating conditions to the new demand. However, if the electrical load is suddenly lost or reduced significantly, a commensurate reduction must be made in the flow of steam through the turbine or the turbine will overspeed, possibly causing turbine damage. The main control system does not possess sufficiently rapid response characteristics to accommodate for such sharp variations in load demand, especially in high power to inertia ratio turbine systems.\nAs is well known, large steam turbines generally include multiple nozzle chambers through which steam is directed into turbine through turbine blades which are rotated thereby. Nozzle chamber activation (i.e., steam admission thereinto) is regulated by valves which open to provide steam flow from steam supply conduits into the nozzle chambers, and close to obstruct steam flow thereinto. Inactive valves are completely closed at valve points in an ideal valve position, while active valves are wide open. In reality, to improve load response, the last valve that has opened is not in the wide open position before the next valve begins to open. A valve point is defined as a state of steam admission in which each active valve is in the completely open, non-obstructing configuration or each inactive valve is in the completely closed, full obstructing configuration. It can be shown that maximum turbine efficiency can be obtained from the use of an infinite number of valve points which, in turn, requires an infinite number of valves.\nOf course, a finite number of valves must be used on steam turbines with that number of valves being dictated by compromises between improved turbine performance and increasing capital cost for increasing numbers of valves. One or more valves control the flow of steam into each nozzle chamber. Nozzle chamber activation refers to the process of increasing steam flow into the nozzle chambers from the time steam flow thereinto (i.e., completely activated) is achieved. Deactivation refers to the process of decreasing steam flow into the nozzle chambers. When multiple valves are used to regulate steam flow into a single nozzle chamber, those valves typically modulate together. Since such valves modulate together, turbine efficiency is actually a maximum when the nozzle chambers are each in the completely activated or completely deactivated state. Heretofore, the nozzle chambers were activated in a predetermined sequence such that once the nozzle chamber was activated during increasing load on the turbine, it was not deactivated until the load on the turbine decreased. One of the few restraints on nozzle chamber activation sequence was that single shock operation was preferred over double or multiple shock operation. That is, it is usually a preferable practice to activate nozzle chambers such that a newly activated nozzle chamber (i.e., after minimum admission) is circumferentially adjacent at least one previously activated nozzle chamber. One illustrative method for admitting steam into a steam turbine is disclosed in U.S. Pat. No. 4,325,670, issued Apr. 20, 1982 to George J. Silvestri, Jr., assigned to the assignee of the present invention, and incorporated herein by reference.\nOne recurring problem encountered by such turbines, however, is known in the art as low cycle thermal fatigue. With many older turbines being relegated to cycling operations such as load following and on-off or \"two-shifting\" operation, the potential for low cycle thermal fatigue is increased significantly. The problem of low cycle thermal fatigue can nevertheless be minimized in newer turbines by placing individual actuators for each valve in the steam chests of the turbines. Older steam chests, such as those used in the mechanical hydraulic (MH), analog electric hydraulic (AEH), and digital electric hydraulic (DEH) turbine control systems, may not have individual valve actuators, nor may they have sufficient space between the valves to accommodate individual valve actuators. This is especially true in those cases where the actuator incorporates springs necessary to insure rapid closure of the valves during turbine trips.\nA solution to such problems would be the wholesale but costly replacement of the steam chests. However, such wholesale replacement would not only be costly, but would also be extremely time consuming, thereby leading to an inordinate amount of turbine down time. Another solution would be to modify the steam chests of such turbines as described in the above referenced U.S. Ser. No. 107,735, filed Oct. 13, 1987. That is, turbine operations may be improved in a conventional steam turbine by improved steam chest means. The steam turbine, as is typical, has a casing including inlet means for receiving a flow of steam by the steam chest means. Comprising a plurality of valves each of which are set for a minimum admission of the flow of steam to the inlet means below 100%, bar lift means for actuating at least one pair of the valves, high pressure means for actuating remaining ones of the plurality of valves, and means for controlling the bar lift means and high pressure means, the steam chest means thereby allows the turbine to be transferred between a full-arc (or maximum) admission mode and a partial-arc (lower level) admission mode.\nIn steam chests of the internal bar lift type, the bar is shortened or removed such that only the two innermost valves of a four-valve steam chest are still actuated by the bar lift means, while the two outboard valves at each end of the steam chest are replaced with ones having individual high pressure actuators. For those steam chests of the end bar or external bar lift type, the pivot on the fixed end of the bar is replaced with another servomotor such that the actuator rod of the new servomotor incorporates the pivot for the external bar. By a combination of lifts of the existing servomotor and the new servomotor, it is possible to operate at full-arc admission at start up and to make the transition from full (or maximum) to partial-arc (and vice versa) at whatever level of load is desired and whatever value of partial-arc admission is consistent with first stage requirements and optimum loading conditions.\nIt is well known that low load and part load operation of steam turbines with sliding throttle pressure not only reduces the above mentioned low cycle thermal fatigue, but also improves the heat rate. In particular, operation in a hybrid mode (i.e., a combined mode of operation with constant pressure-sequential valve and sliding throttle) results in a maximum heat rate benefit while reducing the change in first stage exit temperature, thereby reducing low cycle thermal fatigue. With hybrid operation, a partial-arc admission turbine is operated in the upper load range by activating individual valves to effect load changes along with constant throttle pressure operation. As load is reduced, when a particular valve point is reached, valve position is held constant and throttle pressure is varied or slid to achieve further load reductions. On units with essentially 100% admission at maximum load, hybrid operation with a 50% minimum first stage admission achieves the heat rate benefit of constant throttle pressure operation. Additionally, when valve loop losses are considered, hybrid operation has superior thermal performance to partial-arc designs operating with constant throttle pressure and having admission points below 50% at loads below from 65 to 70& of maximum value. For units with considerably less that 100% admission at maximum load, optimum hybrid operation is achieved at the valve point where half of the valves are wide open and half are closed. It is for this reason that the above referenced U.S. Ser. No. 107,735 provides one method and apparatus for a valving sequence on turbines having steam chests without individual valve actuators in such a manner that the valves which correspond to 50% first stage admission (or half of the total number of valves) all open simultaneously.\nHowever, start up procedures that increase rotor life require a different operating mode than hybrid mode operation. Full-arc admission during turbine roll, for example, has proven beneficial for rotor warmup and more uniform heating as well as reducing the steam-to-metal temperature mismatches that increase low cycle thermal fatigue. It has also been noted that maintaining full-arc admission operation beyond synchronization of the turbine up to some level of load can be beneficial. Full-arc admission operation at part load, however, cannot be achieved on turbines having steam chests without individual actuators for which the valves are set for minimum first stage admissions below 100%. It has also been noted that an expected increase in rotor life is achievable when the transfer from full to partial-arc is made during the loading cycle as compared to full-arc admission operation all the way to full load.\nIt is, therefore, apparent that a steam chest having the capability of valve transfer from full to partial-arc admission and vice versa would be extremely desirable for turbines utilized in cycling operations. As mentioned herein above, U.S. Ser. No. 107,735, to which the present application is a continuation-in-part, discloses various embodiments of one method and apparatus for achieving such results. The present application, on the other hand, discloses an alternative method and apparatus for achieving the same."}
{"text": "The present invention relates to an electronic flash control device for a TTL light control type camera, and more particularly to an electronic flash control device which is free from the difficulty that light emission by the electronic flash is stopped by a noise component occurring before a suitable quantity of output light is obtained.\nWhen it is necessary to photograph an object having a low luminance or to correct the amount of exposure, an electronic flash is generally used to increase the luminance of the object. The electronic flash is either built into the camera or is provided separately and can be connected to the camera. Whether the electronic flash is built into the camera or provided separately from the camera, it is essential to control the light emitting conditions of the electronic flash so that the amount of light applied to the object by the electronic flash has a predetermined value. In order to meet this requirement, the camera is provided with an electronic flash control device.\nThe electronic flash control device will be described with reference to a TTL light control type camera employing an electronic flash. FIG. 4 shows an example of a camera of this type. More specifically, FIG. 4 is a block diagram outlining the arrangement of a camera having an automatic focusing (AF) function and a built-in electronic flash proposed by the present applicant.\nFirst, the arrangement of the camera shown in FIG. 4 will briefly be described. In FIG. 4, reference numeral 31 designates a camera body, and 11 represents a photographic lens mounted on the camera body 31.\nThe photographic lens 11 includes a lens system 15 including a focusing lens 13 which is movable along the optical axis and contributes to the focusing operation, and a drive force transmitting mechanism 17 for transmitting a drive force from a drive source provided for the camera body 31 to the movable lens 13. The photographic lens 11 further includes a lens ROM (read-only memory) 19 storing aperture value data of the photographic lens, position data of the movable lens 13, and data as to whether or not operating the electronic flash built into the camera body is suitable for the photographic lens, and a group of electrical contacts 21 which are connected to electrical contacts 57 on the camera body 31.\nOn the other hand, the camera body 31 includes an optical system having a main mirror 33, an auxiliary mirror 35, a focusing screen 37, and a pentagonal prism 39, an image pickup section 41 used for an automatic focusing operation, a drive mechanism 43 for driving the movable lens 13 in the photographic lens 11, a light detecting element 45 used for automatic exposure (AE) control for the purpose of TTL light control, a light detecting element 47 operated when the electronic flash is used, a central display section 49 for displaying camera conditions, a display unit 51 provided in the view-finder for indicating an AF or AE mode, the built-in electronic flash 53, a sequence motor 55 for winding or rewinding the film, a group of electrical contacts 57 provided on the side of the camera body which are connected to the group of electrical contacts 21 on the side of the photographic lens, a release switch 59, and a synchronizing contact 61 such as an X contact.\nThe camera body 31 further includes a microcomputer, namely, an IPU (indication processing unit) 71 for controlling the central display section 49, a PCU (power control unit) 73 having an E.sup.2 PROM 73a for interfacing with the image pickup section 41 and controlling the sequence motor 55 and an AF motor 43, the lens stop and the shutter release magnet, another microcomputer, namely, a DPU (data processing unit) 75 for performing photometric arithmetic operations and controlling the display unit 51 in the viewfinder, etc., and a central control microcomputer, namely, a CPU (central processing unit) 77. The CPU 77 controls the IPU 71, the PCU 73, the DPU 75, and the lens ROM 19 in the photographic lens 11.\nThe above-described built-in electronic flash 53 is provided above the pentagonal prism 39 and substantially at the middle of the camera body 31 so that the light flash can be readily applied to an object. When the electronic flash 53 is not in use, its light flash emitting section is retracted into the camera body. When it is in use, the light flash emitting section is protruded therefrom by operating a push button (not shown). The camera is designed so that an auxiliary electronic flash can be additionally connected thereto.\nNext, an electronic flash control device for the above-described camera will be described. FIG. 5 is a diagram showing the above-described electronic light flash 53 and the electronic flash control device in detail. In FIG. 5, reference numeral 81 designates the electronic flash control device. The device 81 includes as essential components the CPU 77, the DPU 75, the PCU 73, the light detecting element 47, the release switch 59 and the X contact 61. The electronic flash 53 includes a booster circuit 53a, a main capacitor 53b, a light emitting circuit 53d having a light emitting tube 53c, and a light emission stopping circuit 53e. The electronic flash 53 is connected to the electronic flash control circuit 81 through a trigger signal terminal 53f and a quench signal terminal 53g. When necessary, the aforementioned auxiliary electronic flash is connected to the electronic flash control circuit 81 through auxiliary electronic flash connecting contacts corresponding to the above-described contacts 53f and 53g which are provided at predetermined positions on the camera body 31.\nThe operation of the electronic flash control device 81 thus constructed will be described with reference to the timing diagrams shown in FIGS. 6(A) through 6(J).\nThe release switch 59 is turned on after the camera has selected the electronic flash operation mode. Then, the CPU 77, after carrying out a series of processing operations, starts detecting the variation of the voltage at the X contact. In response to the variation of the voltage at the X contact (SWX in FIG. 6(A)) which occurs when the shutter is fully opened, the CPU 77 supplies a trigger signal X.sub.TRIG shown in FIG. 6(B) to light emitting circuit 53d in the electronic flash 53 via trigger signal terminal 53f. In response to the trigger signal, the light emitting circuit 53d supplies an electrical signal S.sub.X shown in FIG. 6(C) to the light emitting tube 53c. As a result, the light emitting tube 53c emits light. (See FIG. 6(G)).\nOn the other hand, an integrator circuit 75a in the DPU 75 integrates the output signal of the light detecting element 47 in response to an integration start signal S.sub.S shown in FIG. 6(D) which the CPU 77 outputs in response to the trigger signal. The integration voltage S.sub.I of the integrating circuit, as shown in FIG. 6(E), gradually increases with time, and abruptly increases with the emission of light from the light emitting tube. The integration voltage S.sub.I is applied to one input terminal of a comparator 75b.\nThe CPU 77 applies a digital signal S.sub.D to a D/A (digital to analog) converter 75c prior to every photographing operation. The digital signal represents a suitable exposure value determined in accordance with the ISO data of the film loaded in the camera or exposure correction data when available. The digital signal S.sub.D is subjected to digital-to-analog conversion by the D/A converter 75c, the output of which is applied to the other input terminal of the comparator 75b. The output terminal of the comparator 75b is connected to one input terminal of an AND circuit in a quench signal output regulating circuit 75d in the DPU 75.\nWhen, in the circuit thus constructed, the integration voltage S.sub.I reaches a reference value S.sub.O, the comparator 5b applies a signal to the AND circuit in the quench signal output regulating circuit 7.5d so that the quench signal is generated (in a manner to be described in more detail later). When, with an enable signal S.sub.E applied to the other input terminal of the AND circuit by the CPU 7, the integration voltage S.sub.I reaches the reference value, the regulating circuit 75d applies a quench signal Q.sub.U shown in FIG. 6(F) to the light emission stopping circuit 53e. In response to the quench signal, the emission of light of the electronic flash 53 is quickly suppressed as shown in FIG. 6(G).\nThe present applicant has conducted extensive research on the electronic flash control circuit described above and recognized that the output timing of the enable signal S.sub.E is very important. That is, depending on the enable signal, the control circuit may operate erroneously.\nWhen the CPU 77 supplies the enable signal S.sub.E as shown in FIG. 6(H) to the quench signal output regulating circuit 75d before a noise pulse N.sub.T occurs within the electrical signal S.sub.X in response to the trigger signal as shown in FIG. 6(C), the circuit is operated erroneously and a false quench signal Q.sub.UF is produced as shown in FIG. 6(I), so that no light is emitted (See FIG. 6(J).) In practice, such noise pulses are often produced by the light emitting tube, and accordingly, in order to perform the emission of light correctly, it is essential to use a light emitting tube having excellent characteristics. As a result, during manufacture, it is necessary to perform an additional operation of selecting light emitting tubes suitable for the electronic flash control device.\nIn view of the foregoing, an object of the present invention is to provide an electronic flash control device which is free from the difficulty of the emission of the light flash being stopped by a noise component before a suitable amount of light is obtained."}
{"text": "U.S. Pat. No. 4,562,082 describes the .alpha.-factor precursor gene of Saccharomyces cerevisiae and purports to show DNA constructs of the S. cerevisiae leader sequence and heterologous genes that are useful for the secretory expression of the heterologous genes. EPA 83306507.1 describes the actual use of the S. cerevisiae .alpha.-factor leader sequence to obtain secretory expression of human epidermal growth factor in yeast. Copending U.S. patent application Ser. No. 522,909, filed Aug. 12, 1983, which is, in part, commonly owned with the present application, describes the use of \"spacer-less\" S. cerevisiae .alpha.-factor leader sequences to obtain more efficient secretory expression of heterologous genes.\nScience (1985) 229:1219-1224 reports the use of S. cerevisiae .alpha.-factor sequences to direct the secretion of prochymosin. It was reported that most of the prochymosin produced was present in an intracellular insoluble form rather than being secreted. Although the invertase leader was more effective at directing secretion of activatable prochymosin, an extensive mutant screen was necessary to yield reasonable amounts of prochymosin secretion.\nThere has been no previous demonstration of the production of mating pheromones (.alpha.-factor and .alpha.-factor) by Kluyveromyces lactis cells or other Kluyveromyces species. This failure to identify such peptides in K. lactis may, in part, be due to the fact that mating in K. lactis is inducible; whereas in S. cerevisiae the mating pheromones are produced constitutively.\nNotwithstanding the availability of the mentioned leader sequences for directing secretion of heterologous polypeptides in yeast, there continues to be a need for alternative sequences that may afford more efficient or more practical production of particular polypeptides. Accordingly, applicants sought to determine whether a K. lactis .alpha.-factor exists and, if so, whether it might be useful for directing secretion of heterologous polypeptides, such as prochymosin, in yeast."}
{"text": "An ultrasonic diagnosis device sends an ultrasonic signal to an inspection target, and visualizes an internal status of the inspection target based on reflected wave information thereof. This device is used for nondestructive inspection in a medical treatment region or an industrial region.\nAs to the ultrasonic diagnosis device, first, by processing a signal received via a probe, two data (identical phase component data (I data) and orthogonal phase component data (Q data)) are calculated. Then, based on I data and Q data (Hereinafter, they are called “IQ data”), internal data (A-mode data, B-mode data, color Doppler data (including a direction, an averaged flow velocity, a power, a distribution), M-mode data) are calculated. By applying a coloring or a scan transform to the internal data, the inspection target is visualized as image data. Moreover, in a diagnosis using this device, the power Doppler method for visualizing power data is also used. Hereinafter, for convenience' sake, the color Doppler data and the power Doppler data are called “color Doppler data”.\nRegularly, the ultrasonic diagnosis device prepares a plurality of methods for visualizing IQ data and internal data as image data. A user can select any from the plurality of methods. Furthermore, some ultrasonic diagnosis device prepares a function to preserve the image data visualized. As a result, after the diagnosis in real time, the user can confirm the image data visualized in off-line. In order to switch the plurality of methods in off-line, by preserving not the image data already visualized but the reflected wave information (IQ data, or internal data) not visualized yet, the visualization processing needs to be applied to the preserved data again.\nHowever, a data quantity of the reflected wave information is much larger than that of the image data visualized. Accordingly, preservation of all reflected wave information unprocessed is not realistic. Furthermore, if compression processing is applied to all reflected wave information inputted, processing quantity is enormous. As a result, processing time thereof becomes very long."}
{"text": "Conventionally, trucks in a conveyor are driven by a driving chain. In such a conveyor, the driving chain is moving at all times even while the trucks are stored in storage lines, which promotes occurrence of wear of the driving chain as well as noises. Moreover, in order to start the trucks, it is necessary to engage the driving chain with the pressure plate provided to the trucks, which causes a shock at the time of engagement. Such a shock of course becomes greater in proportion to the speed of the driving chain. For these reasons, the upper limit of the speed of the conventional conveyor has been some 18 meters per second. Thus, the conventional chain drive conveyors are not suitable for high speed conveyance.\nOne solution of this problem is to use a conveyor driven by linear motor. In known linear motor drive systems, primary members of linear induction motor are provided on the floor and reaction plates which constitute the secondary members of a linear induction motor are provided to the trucks, respectively; and by maintaining the gap between the two members accurately and energizing the primary members, the trucks are electromagnetically propelled.\nHowever, if there is a horizontal curve in the conveying path, the locus of the reaction plate is likely to be diverted from the primary members; and if there are curves or undulations in the vertical direction, the reaction plates are subjected to movements up and down. Thus, there occurs unevenness of the gap, which will result in instability or reduction in the propelling force. Consequently, conventional linear motor drive systems have been used only in straight conveying lines.\nAccordingly, one object of this invention is to provide a linear motor driven trolley conveyor which is capable of stably driving trucks even where there are horizontal curves or slopes in the conveying path."}
{"text": "1. Field of the Invention\nThe present invention relates to microprocessor architecture in general and in particular to a microprocessor architecture capable of supporting multiple heterogeneous microprocessors.\n2. Description of the Related Art\nA computer system comprising a microprocessor architecture capable of supporting multiple processors typically comprises a memory, a memory system bus comprising data, address and control signal buses, an input/output I/O bus comprising data, address and control signal buses, a plurality of I/O devices and a plurality of microprocessors. The I/O devices may comprise, for example, a direct memory access (DMA) controller-processor, an ethernet chip, and various other I/O devices. The microprocessors may comprise, for example, a plurality of general purpose processors as well as special purpose processors. The processors are coupled to the memory by means of the memory system bus and to the I/O devices by means of the I/O bus.\nTo enable the processors to access the MAU and the I/O devices without conflict, it is necessary to provide a mechanism which assigns a priority to the processors and I/o devices. The priority scheme used may be a fixed priority scheme or a dynamic priority scheme which allows for changing priorities on the fly as system conditions change, or a combination of both schemes. It is also important to provide in such a mechanism a means for providing ready access to the memory and the I/O devices by all processors in a manner which provides for minimum memory and I/O device latency while at the same time providing for cache coherency. For example, repeated use of the system bus to access semaphores which are denied can significantly reduce system bus bandwidth. Separate processors cannot be allowed to read and write the same data unless precautions are taken to avoid problems with cache coherency."}
{"text": "Drinking water has been, and continues to be, heavily treated for bacteria and other microscopic organisms that may cause infection in humans and other animals subsequent to consumption. In order to disinfect water supplies, halogenated materials have been introduced therein that have proven more than adequate for such a purpose. Unfortunately, although such halogenated compounds (chlorinated and chloraminated types, primarily) exhibit excellent disinfection capabilities, when present within aqueous environments at certain pH levels these halogenated compounds may generate byproducts that may themselves create health concerns. The United States Environmental Protection Agency (USEPA) in fact regulates four types of trihalomethanes (THM4) within drinking water. These THM4 are chloroform, bromoform, dibromochloromethane, and bromodichloromethane. Removal of such compounds from drinking water is not possible as for typical chlorinated and brominated disinfecting compounds, at least not at the same reliability level as for the disinfecting agents. Thus, residual amounts may remain within treated water supplies that may require further removal processing to be undertaken, or avoidance of ingestion if necessary. Of course, if the level of contamination is sufficiently low, initiation of such potentially expensive removal steps would be unwise from an economic perspective.\nThe USEPA currently has set a maximum contaminant level for these THM4 in drinking water of 0.080 mg/L. It is thus important to reliably analyze and measure the total amount of such contaminants in order to determine if removal if necessary.\nThe USEPA has instituted its own testing methods for such a purpose. For THM4, the primary test method is 502.2. The USEPA 502.2 method measures for individual and total THM4 as well as other volatile byproducts. This method utilizes a TRACOR® 540 gas chromatography with Hall/PID detectors, a Tracor LSC-2 sample concentrator, and a TEKMAR® 2050 Autosampler. Such a system is, again, very effective at measuring drinking water samples at the source, but remote analyses are not readily available as the entire system is too cumbersome to move to locations along a drinking water line. As such, on-line analysis through these protocols is difficult, expensive, and labor intensive to implement.\nMeasurement at the source (i.e., within a water purification plant location) may be effective for system-wide average readings; however, in the large supplies of water at such locations, the chances of proper sampling to that effect may be suspect since the contaminants may be present in varied locations, rather than definitely mixed throughout the tested water supply itself. Additionally, testing may not uncover the actual level of residual THM4 disinfection byproducts prior to the water supply being disbursed to distant dispense sites (transfer pipes, homes, schools, businesses, etc.). In any event, there is a relatively new rule in place that requires utilities to provide evidence of compliance with trihalomethane levels at multiple locations, rather than a straightforward system-wide average. There is thus a drive to implement remote testing via real-time, on-line methods for water supply THM4 contaminant level measurements.\nSuch a desirable on-line procedure has been difficult to achieve, however, particularly as it pertains to the determination of not only the total amount of THM4 within water supplies, but also the amount of each species of the same THM4 groups present within the tested water source. High performance liquid chromatography, utilizing electrospray ionization-mass spectrometry or ultraviolet absorbance detector, has been attempted. The sensitivity and selectivity of the high performance liquid chromatography methods are easily sacrificed without the cumbersome preparations in place, therefore requiring operator intervention during analysis. Again, this issue leads to serious drawbacks when on-line implementation is attempted.\nAnother methodology that has proven effective to a degree is post-column reaction-ion chromatography. This has shown promise, but only in terms of quantifying bromate ion concentrations in drinking water samples at a single microgram per liter level. This dual selectivity form (separation by ion chromatography column as well as the selective reaction with the post-column reagent with the analyte) offers an advantageous test method over the others noted above, except for the presence of more common anions, specifically chloride, at much higher concentrations within the sampled drinking water supply (mg/L instead of microgram). Separation of the THM4 species from other halogenated compounds (such as haloacetic acids), however has been problematic and caused certain degrees of interference in measuring total levels of both types of compounds within drinking water samples. Despite this problematic limitation, it was determined that fluorescence detection provided a much-improved detection protocol in comparison with ultraviolet and mass spectrometry possibilities, particularly for haloacetic acid concentrations. Thus, although such a fluorescence method of detection, coupled with the post-column reaction (again with nicotinamide reactant) and ion chromatography, exhibited the best results in terms of an on-line test method for haloacetic acid drinking water contaminant measurement levels, there remained a definite need for improvements in individual and total trihalomethane measurements and identifications within such test samples. To date, however, there has not been an analytical test protocol that has permitted implementation of such a system within an on-line real-time monitoring procedure with an acceptable degree of reliability. An automated system that provides such versatility and reliability has simply not been forthcoming within the pertinent art."}
{"text": "The present invention relates to a method for measuring the amount of gas that leaks through sealed packages. More specifically, the invention relates to gas leakage through packages that have been sealed by a cover of porous material fabricated as a mat of polyethylene fibers. This material acts as a permeable membrane to gases, but an impermeable membrane to bacteria. The membrane comprises a layer having pores which provide a tortuous path to the passage of bacteria; the material is commonly sold under the trademark designation xe2x80x9cTYVEC.xe2x80x9d The packages which use this material are typically semi-rigid plastic cases which protect medical devices and appliances after manufacture and before actual use.\nThe invention relates to U.S. Pat. No. 5,939,619, issued Aug. 17, 1999, entitled xe2x80x9cMethod and Apparatus for Detecting Leaks in Packaging,xe2x80x9d and U.S. Pat. No. 6,050,133, issued Apr. 18, 2000, entitled xe2x80x9cMethod and Apparatus for Detecting Leaks in Packaging.xe2x80x9d Both of these patents are owned by the assignee of the present invention. The present invention also relates to co-pending application entitled xe2x80x9cMethod for Measuring Gas Leakage From Sealed Packages,xe2x80x9d Ser. No. 09/676,621, filed Oct. 2, 2000, and owned by the assignee of the present invention. The present invention permits a measurement of leakage that is wholly non-destructive to the package.\nGas sterilization is widely used for medical devices that must be sterile at the time of use, but cannot be subjected to sterilization by the application of high temperatures. Examples of such medical devices include cardiac pacemakers and catheter-based monitoring devices such as blood pressure probes. Typically, the medical device is sealed within a package that is permeable to gases but impermeable to bacteria. The package is then placed in a gas sterilization chamber, and a sterilizing gas such as ethylene oxide is introduced into the gas-permeable package to achieve sterilization. The sterilizing gas is then removed from the package, leaving the interior of the package sterile and non-toxic.\nIn a typical design, the medical device is placed within a thermoformed rigid plastic tray equipped with a flat sealing flange. A sheet of gas-permeable membrane, such as DuPont TYVEK(copyright) 1073-B (medical grade) brand membrane, which is available from E. I. duPont de Nemours and Co., is then sealed to the sealing flange, typically by using an adhesive. The integrity of the seal is critically important to maintaining sterility. Leaks can result from incorrect setting of parameters in the automated sealing process, or from physical defects such as burrs on the face of the sealing equipment.\nAccording to the known practice described in the prior art patents listed herein, a temporary barrier is formed over the gas-permeable layer, wherein the temporary barrier has an aperture with the gas-permeable layer to temporarily seal the gas-permeable layer except where the aperture is located. A tracer gas is applied under low pressure through the aperture so that it can enter into the interior chamber of the package. The entire package is placed into a larger sealed second chamber, and the concentration of tracer gas in the second chamber is measured, outside the package, to thereby measure the amount of tracer gas which has leaked through the package, presumably via leaks in the sealing flange, although leakage can also occur through pinhole defects in the plastic tray itself.\nThe methods described in the foregoing patents provide very accurate measurements and evaluations for sealed packages, under controlled conditions. However, in many applications, it is not necessary to achieve a high degree of accuracy in the leakage measurement, but is desirable to provide a quick evaluation of leakage as a production line test, to determine whether packages are leaking excessively. Excessive leakage is defined in terms of leakage beyond a predetermined range of acceptability, as a pass/fail parameter, and the precise degree of leakage does not need to by quantified. Methods for making this type of determination should produce results more quickly and at a lesser cost.\nThe present invention provides a pass/fail test which can be quickly performed at considerably less cost in terms of test equipment and testing time, because it relies on measurement of internal package gas pressure, and specifically pressure drop, caused by leakage of gas from within the package. The method of the present invention can be performed in a short time, perhaps 30-60 seconds, using very much less expensive equipment than prior art methods.\nThe present invention comprises a method for testing sealed packages of the type described herein for leakage, according to a more efficient and faster process than is known in the prior art. The method requires a barrier layer to be overlaid on the gas-permeable membrane of the package, the barrier layer having a single aperture to the gas-permeable membrane. A pressure-sensitive monitor is affixed over the aperture, to monitor the pressure inside the package, and more particularly, to monitor the relative drop in pressure caused by leakage from the package. The internal package gas pressure is initially incrementally increased as a step of the method, and the incremental drop in pressure is measured as a function of time, usually by determining whether pressure drops incrementally between two desired set points within a predetermined time increment.\nIt is a principal object and advantage of the present invention to provide a method for measuring leakage through packages under faster test conditions than have heretofore been known.\nIt is another object and advantage of the present invention to provide a leakage test method which does not destroy the package being tested.\nOther and further objects and advantages of the invention will become apparent from the following specification and claims and with reference to the appended drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a MIS device, and particularly to an improvement of so-called LDD (Lightly Doped Drain) structure of a PMOS transistor and a NMOS transistor formed on the same substrate and a method of manufacturing the improved LDD structure.\n2. Description of the Background Art\nA device provided with multi-layered structure of metal-insulator-semiconductor is generally referred to as an MIS device. An MOS device is a kind of MIS device, in which an oxide film is used for an insulator of the MIS device.\nSemiconductor devices in which a p-channel MOS field effect transistor (referred to as a pMOS transistor hereinafter) and a n-channel MOS field effect transistor (referred to as a NMOS transistor hereinafter) are formed on a single semiconductor substrate include a CMOS semiconductor device, for example. FIG. 6 is a sectional structure view of a conventional CMOS semiconductor device. A p well region 8a and a n well region 8b in contact with each other are formed on a surface of a silicon substrate 6. Thick field oxide films 7 are formed at given positions on the main surface of the silicon substrate 6. The p well region 8a surface surrounded by the field oxide film 7 forms a n MOS transistor forming region, and the n well region 8b surface surrounded by the field oxide film 7 forms a pMOS transistor forming region.\nThe nMOS transistor 20 has a gate electrode 1a on the surface of the p well region 8a with a thin gate insulating film 9 disposed therebetween. Sidewall spacers 3a and 3a composed of insulating films are formed on sidewalls of the gate electrode 1a. Also, a pair of source/drain regions including n.sup.- impurity regions 2a, 2a of relatively low concentration and n.sup.+ impurity regions 5a, 5a of relatively high concentration are formed in the p well region 8a surface. Such structure of source/drain regions is referred to as LDD structure. Source/drain interconnection layers 11, 11 are respectively connected to the surfaces of the n.sup.+ impurity regions 5a, 5a of the source/drain region.\nAlso, the pMOS transistor 30 includes a gate electrode 1b on the surface of the n well region 8b with a thin gate insulating film 9 disposed therebetween. Sidewall spacers 3b, 3b formed of insulating films are formed on the sidewalls of the gate electrode 1b. Furthermore, source/drain region composed of p.sup.- impurity regions 2b, 2b of relatively low concentration and p.sup.+ impurity regions 5b, 5b of relatively high concentration are formed at the surface of the n well region 8b. This source/drain region forms the so-called LDD structure. Moreover, source/drain interconnection layers 11, 11 are respectively connected to the p.sup.+ impurity regions 5b, 5b of the source/drain region. The silicon substrate 6 surface in which a transistor is formed is covered with a thick interlayer insulating layer 10.\nNext, a method of manufacturing the above-described CMOS semiconductor device will be described referring to FIGS. 7A-7H.\nFirst, referring to FIG. 7A, a p well region 8a and an n well region 8b are formed on a silicon substrate 6 surface, and thick field oxide films 7 are formed at given regions on the silicon substrate 6 surface. Furthermore, gate insulating films 9 are formed on the surfaces of the p well region 8a and the n well region 8b. A polysilicon layer 12 is deposited on the surface of the gate insulating film 9.\nNext, referring to FIG. 7B, the polysilicon layer 12 is patterned into a rectangular form using the photolithography method and the anisotropic etching method to form the gate electrode 1a of the nMOS transistor and the gate electrode 1b of the pMOS transistor.\nMoreover, referring to FIG. 7C, after covering the surface of the p well region 8a with a resist pattern 4a, p type impurity ions 15 such as boron (B) are implanted into the n well region 8b using the gate electrode 1b as a mask. By this ion implantation, p.sup.- impurity regions 2b, 2b of low concentration are formed on the n well region 8b surface.\nFurthermore, as shown in FIG. 7D, after removing the resist pattern 4a on the p well region 8a surface, a resist pattern 4b is now formed covering the surface of the n well region 8b. Subsequently, n type impurity ions 16 such as phosphorus (P) or arsenic (As) are directed on the p well region 8a surface using the gate electrode 1a as a mask to form n.sup.- impurity regions 2a, 2a of low concentration.\nFurthermore, as shown in FIG. 7E, after removing the resist pattern 4b, an oxide film 13 is deposited all over the surface of the silicon substrate 6 using the low pressure CVD (Chemical Vapor Deposition) method.\nSubsequently, referring to FIG. 7F, the oxide film 13 is anisotropically etched to form sidewall spacers 3a, 3b having the same film thickness on the sidewalls of the gate electrodes 1a and 1b.\nNext, referring to FIG. 7G, after forming a resist pattern 4c covering the surface of the n well region 8b again, n type impurity ions 17 such as arsenic are implanted using the gate electrode 1a and the sidewall spacers 3a as masks. Subsequently, after removing the resist pattern 4c, thermal treatment is applied for activating the implanted ions. Thus, n.sup.+ impurity regions 5a, 5a of high concentration are formed in the p well region 8a surface. The LDD structure of the source/drain region of a nMOS transistor is then completed.\nFurthermore, referring to FIG. 7H, a new resist pattern 4d is formed covering the surface of the p well region 8a. P type impurity ions 18 are directed on the n well region 8b surface using the gate electrode 1b and the sidewall spacers 3b as masks. Next, after removing the resist pattern 4d, applying thermal treatment, the ions implanted into the n well region 8b are activated. In this way, the p.sup.+ impurity regions 5b, 5b are formed. Thus, the LDD structure of source/drain region of a pMOS transistor is completed in the above steps.\nSubsequently, an interlayer insulating film 10 is formed on the surface of the silicon substrate 6, and a contact hole is formed at a given position. An interconnection layer 11 is then formed at a given position through the contact hole to finish the process of manufacturing a CMOS semiconductor device (not shown).\nThe progress of fine processing technique of element structure is pointed out as technical background of a CMOS semiconductor device with such LDD structure as described above. The tendency of element structure miniaturization is seen in the aspects such as shortening of the gate length according to the scale down rule and forming of shallow junction region in a MOS transistor. The \"scale down rule\" is described in VLSI Electronics Microstructure Science, Volume 18, \"Advanced MOS Device Physics\", Academic Press. Inc. 1989. Also, it means shortening of the gate length of a MOS transistor, that is, shortening of a channel, which causes various problems so-called short channel effect. That is to say, for example, hot carriers are generated due to the high electric field produced in the vicinity of a drain in a short channel MOS transistor, and a portion of the hot carriers are captured in a trap in the gate insulating film to form a new level. As a result, characteristic degradation such as shift of a threshold value voltage of a MOS transistor and a decrease of mutual transconductance are caused. The high electric field produced in the vicinity of a drain also caused degradation of the drain breakdown voltage due to the avalanche breakdown. The LDD structure of a MOS transistor is a device proposed to solve such problems. Especially, it restrains generation of the high electric field by forming an impurity region with moderately changing concentration distribution in the vicinity of the drain to obtain high avalanche breakdown voltage and a decrease of reliability degradation due to hot carriers.\nIn a miniaturization method of CMOS semiconductor devices, it is needed to reduce the gate length L.sub.Gp of pMOS transistor 30 and the gate length L.sub.Gn of nMOS transistor 20 and make the lengths equal. As the gate lengths L.sub.Gp, L.sub.Gn of pMOS transistor 30 and nMOS transistor 20 get longer, the gate capacitance increases to increase the RC constant, with the result that operation of transistors are delayed. Accordingly, the channel lengths Ln and Lp defined in a self aligning manner utilizing gate electrodes 1b and 1a are also preferably reduced to have the equal lengths.\nHowever, turning to FIG. 6 again, the pMOS transistor 30 and the nMOS transistor 20 are provided with low concentration impurity regions 2a, 2b of the LDD structure having different shapes in a CMOS semiconductor device manufactured through the manufacturing steps described-above. That is to say, in the pMOS transistor 30, the diffusion length of the low concentration p.sup.- impurity region 2b is extremely small as compared to that of the nMOS transistor 20. Also, the channel length L.sub.p between the high concentration p.sup.+ impurity regions 5b, 5b are shorter than the channel length L.sub.n of the nMOS transistor 20. This is because, boron (B), impurity forming the source/drain region of the pMOS transistor 30, has a larger diffusion coefficient as compared to that of phosphorus and arsenic forming the source/drain region of the nMOS transistor 20. This is understood by comparing the steps shown in the above FIGS. 7G and 7H. In other words, in the nMOS transistor 20, the high concentration n.sup.+ impurity regions 5a, 5a formed in a self-align manner with respect to the sidewall spacers 3a slightly diffuse under the sidewall spacers 3a by the thermal treatment for activation. On the other hand, in the pMOS transistor, the high concentration p.sup.+ impurity regions 5b, 5b formed in a self-align manner with respect to the sidewall spacers 3b widely diffuse under the sidewall spacers 3b in the thermal treatment for activating. In this way, the diffused high concentration p.sup.+ impurity regions 5b, 5b cover the low concentration p.sup.- impurity regions 2b, 2b regions, so that the diffusion length of the low concentration p.sup.- impurity region 2b is decreased. Thus, punch through phenomena easily occurs between the pair of high concentration p.sup.+ impurity regions 5b, 5b narrowed due to diffusion of impurity of large diffusion coefficient. This punch through phenomena is further facilitated as the gate length becomes shorter while the device is refined.\nAs described above, the LDD structure of the pMOS transistor and the nMOS transistor composed of impurities of different diffusion coefficients are manufactured employing the sidewall spacers 3a, 3b having the same film thickness. Accordingly, if the film thickness of the sidewall spacers 3a, 3b is selected to be suitable for the LDD structure of the nMOS transistor, for example, resistance to the punch through between the high concentration p.sup.+ impurity regions 5b, 5b is easily degraded in a pMOS transistor. On the other hand, if the sidewall spacers are formed with film thickness suitable for the LDD structure of the pMOS transistor, the drain current between the source and the drain is degraded in the nMOS transistor 20. This is because, as the film thickness of the sidewall spacers 3a is increased, the diffusion length of the low concentration n.sup.- impurity regions 2a, 2a is increased, and the low concentration n.sup.- impurity diffusion layers 2a act as parasitic resistance between the source and the drain to degrade the drain current. In this way, in a conventional CMOS semiconductor device, the sidewall spacers 3a, 3b defining the low concentration impurity regions of the LDD structure are formed with the same film thickness in the pMOS transistor and the nMOS transistor, so that a LDD structure which can satisfy requirements of both of the pMOS transistor 30 and the nMOS transistor 20 could not be obtained. Conventionally, by the sacrifice of source-drain current degradation of the nMOS transistor 20 to some extent to prevent the punch through phenomena of the pMOS transistor 30, the film thickness of the sidewall spacers 3a, 3b has been determined focusing on the operational reliability."}
{"text": "Femtocell is a technology that allows users of wireless devices, such as cell phones, PDAs, laptops equipped with cellular communications capabilities, or the like, to install small base stations in the indoor environment using licensed spectrum in order to increase coverage and/or throughput. Femtocells may provide access between a wireless device and a mobile network operator (MNO) core network. For example, the femtocell may be coupled to the operator core network through the Internet. Femtocell base stations may be installed by customers of MNOs on the customer's premises. Aspects such as energy consumption and electromagnetic emission may therefore be important factors for customer acceptance of femtocell products. FIG. 1 illustrates a femtocell base station 110 operatively coupled to an operator core network 120 via an Internet connection 130. The femtocell 110 may operate as a base station for a wireless device 140. Alternatively, a macrocell 150 may operate as a base station for the wireless device 140.\nSome technical standards related to Femtocell technology include 3GPP TS 25.367 V9.2.0, “Technical Specification Group Radio Access Network; Mobility Procedures for Home NodeB; Overall Description,” 3rd Generation Partnership Project, Stage 2 (Release 9), December, 2009; 3GPP TS 25.467 V9.1.0, “Technical Specification Group Radio Access Network; UTRAN architecture for 3G Home Node B (HNB),” 3rd Generation Partnership Project, Stage 2 (Release 9), December, 2009; and 3GPP TS 25.469 V9.0.0, “Technical Specification 3rd Generation Partnership Project; Technical Specification Group Radio Access Network; UTRAN Iuh interface Home Node B Application Part (HNBAP) signalling,” 3rd Generation Partnership Project, (Release 9), December, 2009.\nThere are several differences between a femtocell base station and a traditional cellular base station. A femtocell base station typically connects to the core network using a broadband Internet connection such as DSL, cable or Optical Fibre. Additionally, a femtocell base station may be placed in a user's home or an office building, and typically does not belong to nor is it under direct control of a MNO. For example, an owner of a femtocell base station may want to restrict access to their femtocell since they pay the expense of maintenance and the broadband connection to their premises. As such, a closed subscriber group (CSG) is defined in 3GPP and other standard bodies (e.g., IEEE 802.16m) for a femtocell base station. In other words, a femtocell base station is typically accessible only to terminals which are members of this base station, except for emergency services.\nThe traditional practice of having radio equipment continually turned on can result in energy being wasted. In addition to wasting energy, always-on radio equipment can increase radio interference and radio emissions in nearby areas.\nSome power-saving techniques have been proposed for non-femtocell base stations, such as macrocells. For example, Ericsson™ has introduced a Base Transceiver Station Power savings feature, wherein, during periods of low network traffic, those parts of the network that are not being used may be powered down. Alcatel-Lucent™ offers a feature called Dynamic Power Save which enables a GSM base station to monitor traffic activity on the network and switch off the amplifier when there is no information to transmit. However, these approaches are unsuitable for femtocell base stations, for example since they may require knowledge of the network traffic load and full ownership of the base stations by the MNO. Furthermore, such solutions may not fully leverage the potential capabilities of a femtocell, for example as enabled by the availability of closed subscriber groups.\nTherefore there is a need for a solution that overcomes at least one of the deficiencies in the art.\nThis background information is provided to reveal information believed by the applicant to be of possible relevance to the present invention. No admission is necessarily intended, nor should be construed, that any of the preceding information constitutes prior art against the present invention."}
{"text": "Electronic devices, including portable electronic devices, have gained widespread use and may provide a variety of functions including, for example, telephonic, electronic messaging and other personal information manager (PIM) application functions. Portable electronic devices comprise several types of devices including mobile stations such as simple cellular telephones, smart telephones, Personal Digital Assistants (PDAs), tablet computers, and laptop computers, that may have wireless network communications or near-field communications connectivity such as Bluetooth® capabilities. In addition, electronic devices are also widely used in personal entertainment and infotainment systems, for example, portable media players and automobile infotainment systems.\nThe popularity of electronic devices is driven by user experiences and the interaction between people and the devices via user interfaces. User Interfaces (UIs) that are user friendly and intuitive, functional and stylish, vivid and life-like drive the attractiveness of the device to a consumer.\nImprovements in the method of generating and presenting user interfaces are desirable."}
{"text": "Composite resins that are typical examples of dental compositions usually contain polymerizable monomer compositions including polymerizable monomers and additives such as a filler, a polymerization initiator, a polymerization inhibitor and a dye. In a composite resin including such components, a filler usually has the largest weight fraction followed by a polymerizable monomer composition and these two components represent a major proportion of the weight of the composite resin. The polymerizable monomer composition serves as a binder for the filler. The properties of the monomers, and the properties of cured products of the monomers are significantly influential on the properties and performance of the composite resin containing the monomer composition, and cured products thereof.\nFrom the points of view of aspects such as the biological safety of monomers and the mechanical strength and wear resistance of cured products, the polymerizable monomer compositions frequently include radically polymerizable polyfunctional methacrylates. Typically, the polyfunctional methacrylate compositions are based on 2,2-bis[4-(3-methacryloyloxy-2-hydroxypropoxy)phenyl]propane (hereinafter, written as Bis-GMA) or 2,2,4-trimethylhexamethylenebis(2-carbamoyloxyethyl) dimethacrylate (hereinafter, written as UDMA), and contain triethylene glycol dimethacrylate (hereinafter, written as TEGDMA) to control the viscosity.\nIn the dental clinical practice, the use of composite resins in the restoration of tooth defects has a long history and is still expanding. However, the mechanical properties of cured composite resins are still insufficient. In particular, the poor strength obstructs the application of the resins to sites subjected to a high stress, for example, the use as molar tooth crowning materials.\nIn recent years, clinical experts strongly demand the expansion of the use of composite resins to such high-stress sites. Therefore, the development of composite resins having higher mechanical properties is an urgent necessity. As mentioned above, the properties of cured products of polymerizable monomer compositions used for composite resins significantly affect the properties of cured products of the composite resins containing the compositions.\nTechniques have been reported in which Bis-GMA and UDMA that are widely used as main components of polymerizable monomer compositions are replaced by other monomers so as to enhance the mechanical strength of cured products of composite resins (Patent Literature 1 and Patent Literature 2).\nFurther, techniques aiming to improve main component monomers have been reported. For example, main component monomers are improved so as to enhance the refractive index of cured products of polymerizable monomer compositions (Patent Literature 3), and main component monomers are improved so as to enhance the degree of polymerization shrinkage between before and after the curing of polymerizable monomer compositions (Patent Literature 4)."}
{"text": "1. Field of the Invention\nThe present invention relates to a tire tread monitoring system and more particularly pertains to a system for detecting the initial signs of tire tread separation while the vehicle is in motion.\n2. Description of the Prior Art\nThe use of tire monitoring and warning devices are known in the prior art. By way of example, U.S. Pat. No. 3,489,998 to O'Neal discloses a pressure warning apparatus for a pneumatic tire. U.S. Pat. No. 3,665,388 to Johnsen illustrates a tire air pressure detection and signalling device. U.S. Pat. No. 3,934,223 to Barabino discloses a tire pressure warning system. U.S. Pat. No. 4,089,225 to Kraska discloses a system for residual tire life prediction. U.S. Pat. No. 4,570,152 to Melon discloses a magnetic tire monitor system. U.S. Pat. No. 4,717,905 to Morrison Jr. discloses a warning system for monitoring a condition such as air. U.S. Pat. No. 4,750,584 to Tanaka discloses a distance measuring device. U.S. Pat. No. 4,814,744 to Collins discloses a low pressure warning system. U.S. Pat. No. 5,436,612 to Aduddell discloses a audible vehicle monitoring apparatus. U.S. Pat. No. 5,793,285 to Wehinger discloses a method and device for monitoring a tire on a vehicle.\nIn this respect, the monitoring system of the present invention substantially departs from the conventional concepts and designs of the prior art, by providing a means specifically designed to detect the initial signs of tire tread separation while the vehicle is in motion."}
{"text": "The invention relates to an electrical appliance having a battery holder, which battery holder has a housing which can accommodate at least one battery and which battery holder can be inserted into a housing of the electrical appliance, locking means being provided for locking the battery holder in the appliance.\nSuch removable battery holders, also referred to as power packs, are used, for example, in electrically chargeable appliances such as drills. An example of this is described in DE-A1-39 02 442. The construction for locking such battery holders in the electrical appliance often comprises a large number of parts, which are usually made of a plastic.\nIt is an object of the invention to provide a removable battery holder for an electrical appliance, which holder has a comparatively small number of parts for locking the battery holder and which has a locking mechanism which is reliable, also in the long run.\nTo this end, the appliance in accordance with the invention is characterized in that the locking means are formed by two spaced-apart resilient plastic fingers which each have a projection and by at least one boss for cooperation with the projections of the resilient fingers which fingers and boss have been provided, respectively, on the housing of the battery holder and on the housing of the appliance or the other way around, and which resilient plastic fingers are supported, in the direction of their resilient action, by resilient limbs of a material having a stiffness greater than the stiffness of the resilient plastic fingers. The plastic fingers and the boss can simply be molded integrally with the housing of the battery holder and with the housing of the appliance, respectively. When the battery holder is fitted into the appliance the boss slides between the resilient fingers, upon which the projections are moved apart and subsequently come together after passage of the boss. The resilient limbs whose stiffness is greater than that of the fingers back up the resilient action of the plastic fingers, as a result of which a reliable operation and locking of the battery holder is achieved. To remove the battery holder from the appliance no release button as in DE-A1-39 02 442 is needed. Moving plastic parts are frequently subject to creep. Creep gives rise to a reduction of the force with which the battery holder is locked in the appliance. The resilient limbs prevent the locking mechanism from being adversely affected by creep. The resilient limbs form an additional part but they prolong the life of the mechanism significantly. Moreover, the resilient limbs have an additional advantage, as will be explained hereinafter.\nPreferably, the material of the resilient limbs is an electrically conductive metal.\nA preferred embodiment is characterized in that the housing on which the boss has been provided has electrically conductive limbs which upon insertion of the battery holder into the electrical appliance automatically make contact with the resilient metal limbs which serve for supporting the resilient plastic fingers. The metal limbs make electrical contact with the terminals of the battery (batteries). In their turn, the electrically conductive limbs are connected to a part to be energized of the electrical appliance."}
{"text": "1. Technical Field\nThis invention relates generally to flexible strapping mechanisms and, more particularly, to a flexible support device for removably securing an ice-pack to various joints and muscles of a human body.\n2. Discussion\nIce-packs or hot-water containers are often used by athletes and non-athletes to reduce swelling and remove soreness from various joints or other areas of the human body, or such swelling or soreness as the result of muscular injury, or muscular inflammation or even arthritis. Securing an ice-pack or hot-water container to a traumatized joint for various periods of time often helps to alleviate the discomfort, swelling and/or soreness associated with such muscular injuries or other degenerative conditions.\nWhen an ice-pack or hot-water container is intended to be applied to a particular joint or muscle for various periods of time, it becomes necessary to incorporate some means for securing the ice-pack or hot-water container securely against the inflamed or sore joint or muscle. Heretofore, this has commonly been accomplished by various types of tape, strapping and bandages which are adapted to be wrapped around the ice-pack or hot-water container and areas of the affected joint or muscle to thereby hold the ice-pack or hot-water container securely against the joint or muscle. Typically, such devices have proved to be time consuming and/or inconvenient to secure to the various joints of the human body. This becomes a particular drawback with athletes who incur minor injuries during sporting events. In such cases ice-packs are often periodically applied to the injured joint of an athlete during intermissions, half-times, etc., when an athlete is temporarily resting. When the athlete must periodically enter and exit a game or other event, the tape, bandage, etc., securing an ice-pack to an injured joint or muscle often has to be quickly removed or applied, while causing only a minimum of discomfort to the athlete.\nPresent devices for the above-described purpose often impede the quick and/or periodic application or removal of an ice-pack. Bandages are often wrapped around a limb or joint several times, with multiple bandages often being required to secure an ice-pack to various parts of the human body such as the shoulders. Such devices are therefore of limited value to athletes, in particular, where ice-packs, hot-water containers, etc., must be applied and removed periodically in a quick and efficient manner.\nIt is therefore a principal object of the present invention to provide a universal ice-pack holder operable to quickly and efficiently secure an ice-pack, hot-water container, etc., to a joint or muscle area of the human body.\nIt is a further object of the present invention to provide a universal ice-pack holder which may be used to secure an ice-pack, hot-water container, etc., to virtually any joint or muscle area of the human body in a quick and efficient manner.\nIt is still a further object of the present invention to provide a universal ice-pack holder which is comfortable to wear in contact with the skin for prolonged periods of time while reducing the risk of rashes or other forms of skin irritation.\nIt is yet another object of the present invention to provide a universal ice-pack holder for securing an ice-pack, hot-water container, etc., to a joint or muscle area of the human body, which is of a relatively simple and inexpensive construction.\nIt is still a further object of the present invention to provide a universal ice-pack holder which is manufactured of materials which stretch and generally conform to the shapes of various limbs of a human body, thereby enhancing the comfort of the ice-pack holder."}
{"text": "The present invention relates to polymer compositions having a specified pH to provide improved dispensing for wrinkle removal and/or reduction. The specified pH prevents staining of treated fabrics and methods for treating fabrics are provided in order to improve various properties of fabrics, in particular, reduction, removal, or prevention of unwanted wrinkles.\nFor preferred polymer compositions containing additional components it is particularly important to maintain the specified pH of the polymer compositions to maintain acceptable dispensing while also preventing precipitate formation during processing.\nWrinkles in textile fabrics are caused by the bending and creasing of the textile material which places an external portion of a filament in a yarn under tension while the internal portion of that filament in the yarn is placed under compression. Particularly with cotton fabrics, the hydrogen bonding that occurs between the cellulose molecules contributes to keeping wrinkles in place. The wrinkling of fabric, in particular clothing and certain household fabrics, is therefore subject to the inherent tensional elastic deformation and recovery properties of the fibers which constitute the yarn and fabrics.\nIn the modern world, with the increase of hustle and bustle and travel, there is a demand for a quick fix which will help to diminish the labor involved in home laundering and/or the cost and time involved in dry cleaning or commercial laundering. This has brought additional pressure to bear on textile technologists to produce a product that will sufficiently reduce wrinkles in fabrics, especially clothing and household fabrics, and to produce a good appearance through a simple, convenient application of a product.\nU.S. Pat. No. 5,573,695, issued Nov. 12, 1996 to E. F. Targosz discloses an aqueous wrinkle removal composition containing a vegetable oil based cationic quaternary ammonium surfactant, and an anionic fluorosurfactant. Similarly, U.S. Pat. No. 4,661,268, issued Apr. 28, 1987 to Jacobson et al. discloses a wrinkle removal spray comprising an aqueous alcoholic composition containing a dialkyl quaternary ammonium salt and a silicone surfactant and/or a fluoro surfactant. U.S. Pat. No. 5,100,566, issued Mar. 31, 1992 to Agbomeirele et al., discloses a method of reducing wrinkles in fabric by spraying the fabric with an aqueous alcoholic solution of an anionic siliconate alkali metal salt. U.S. Pat. No. 4,806,254, issued Feb. 21, 1989 to J. A. Church discloses fabric wrinkle removal aqueous alcoholic solution containing glycerine and a nonionic surfactant. U.S. Pat. No. 5,532,023, issued Jul. 2, 1996 to Vogel, Wahl, Cappel and Ward discloses aqueous wrinkle control compositions containing silicone and film forming polymer.\nIn the present invention, wrinkle control in fabrics, including clothing, dry cleanables, linens, bed clothes, draperies, window curtains, shower curtains, table linens, and the like, is acheived without the need for ironing. The present invention can be used on wet, damp, or dry clothing to relax wrinkles and give clothes a ready to wear or use look that is demanded by today\"\"s fast paced world. The present invention also essentially eliminates the need for touch up ironing usually associated with closet, drawer, and suitcase storage of garments.\nAn additional benefit of using polymer-based compositions in the present invention is that polymers provide improved benefits including any or all of the benefits named in the following list: garment shape, body, rewrinkle prevention, and/or crispness.\nWhen ironing is desired however, preferred compositions of the present invention can also act as an excellent ironing aid. The present invention makes the task of ironing easier and faster by plasticizing fabric fibers and thus making it easier to work wrinkles out of the fabric. When used as an ironing aid, the compositions of the present invention help produce a crisp, smooth appearance, but also retaining a quality of softness.\nThe present invention relates to aqueous wrinkle reducing, removing and/or controlling compositions comprising polymer containing carboxylic acid moieties, that are preferably stable, well-dispersed opaque, translucent, or clear suspensions, dispersions, or solutions with the dispersed or solubilized polymer particulates being very small in particle size, that distribute evenly from dispensers to prevent staining. Specified pH solutions are acceptable if these have the low viscosity that is necessary to provide acceptable dispensing. The present invention also relates to preferred compositions containing, in addition to the essential carboxylic acid containing polymer and carrier, optional, but preferred ingredients, e.g. polyalkylene oxide polysiloxane, fabric care polysaccharides, odor control components, solvent, and minors such as perfume and preservative, adjusted to a specified pH to provide both good dispensing properties and improved stability to shear forces (e.g. stirring during processing or shaking that occurs during transit). The present invention further relates to methods of formulating such compositions, as well as fabric wrinkle control methods and articles of manufacture that comprise such fabric wrinkle controlling compositions. The fabric wrinkle control compositions typically comprise:\n(A) an effective amount to control wrinkles in fabric of a polymer preferably selected from the group of polymers comprising carboxylic acid moieties that can be suspended, dispersed or solubilized at a specified pH range to produce a solution with a viscosity lower than the viscosity of that polymer composition at a pH above the specified pH range and with the viscosity of the solution preferably below about 20 centipoise (xe2x80x9ccPxe2x80x9d), more preferably below about 15 cP, even more preferably below about 12 cP, even more preferably below about 10 cP, still more preferably below about 7 cP and most preferably below about 3 cP, with the polymer incorporated at a level that is at least about 0.001%, preferably at least about 0.01%, and more preferably at least about 0.05%, and still more preferrably at least about 0.1% and even more preferably at least about 0.25% and most preferrably at least about 0.5% and at a level of no greater than about 25%, more preferably no greater than about 10%, even more preferably no greater than about 7%, and still more preferably no greater than about 5% by weight of the usage composition; mixtures of polymers are also acceptable in the present composition; and\n(B) carrier, preferably water.\nThe preferred polymer compositions of the present invention can optionally further comprise:\n(A) optionally, but preferably, silicone compounds and/or emulsions. Silicone compounds that impart lubricity and softness are highly preferred. Silicones that reduce surface tension are also highly preferred. A preferred class of silicone materials includes silicones modified with alkylene oxide moieties compounds;\nmixtures of silicones that provide desired benefits are also acceptable in the present composition;\n(B) optionally, an effective amount of a supplemental wrinkle control agent selected from the group consisting of (1) adjunct polymer (2) fabric care polysaccharides, (3) lithium salts, (4) fiber fabric lubricants, and (5) mixtures thereof;\n(C) optionally, an effective amount of a supplemental surface tension control agent;\n(D) optionally, an effective amount to soften fibers and/or polymer of hydrophilic plasticizer wrinkle control agent;\n(E) optionally, but preferably, at least an effective amount to absorb or reduce malodor, of odor control agent;\n(F) optionally, but preferably, an effective amount to provide olfactory effects of perfume;\n(G) optionally, an effective amount of solubilized, water-soluble, antimicrobial preservative, preferably from about 0.0001% to about 0.5%, more preferably fromabout 0.0002% to about 0.2%, most preferably from about 0.0003% to about 0.1%, by weight of the composition;\n(H) optionally, an effective amount to adjust and control pH of a pH adjustment system;\n(I) optionally, other ingredients such as adjunct odor-controlling materials, chelating agents, viscosity control agents, additional antistatic agents if more static control is desired, insect and moth repelling agents, colorants, whiteness preservative; anti-clogging agents; and;\n(J) mixtures of optional components (A) through (I).\nThe present compositions are preferably essentially free of materials that would soil or stain fabric under usage conditions, or preferably free of materials at a level that would soil or stain fabrics unacceptably under usage conditions.\nThe present invention also relates to concentrated compositions, including liquid, fluid and solid forms of concentrated compositions that may be diluted to form compositions with the usage concentrations for use under usage conditions. It is preferred that the concentrated compositions be delivered in forms that rapidly and smoothly dissolve or disperse to the usage concentration.\nThe present invention further relates to a method of making the present compositions.\nThe present invention also relates to articles of manufacture comprising the present compositions incorporated into a container, such as a spray dispenser, that can facilitate treatment of articles and/or surfaces with said compositions containing wrinkle control agent and other optional ingredients at a level that is effective, yet is not discernible when dried on the surfaces. The spray dispenser comprises manually activated and non-manual powered (operated) spray means and a container containing the wrinkle controlling composition.\nThe present invention also relates to the combining the composition with a substrate and/or device capable of containing said composition for release at a desirable time in a fabric treatment process to create an article of manufacture. Such articles of manufacture can facilitate treatment of fabric articles and/or surfaces with said pH adjusted polymer compositions containing wrinkle control agent and other optional ingredients at a level that is effective, yet not discernible when dried on the surfaces of said fabric. The article of manufacture can operate in mechanical devices designed to alter the physical properties of articles and/or surfaces such as, but not limited to, a clothes dryer or mechanical devices designed to spray fabric care compositions on fabrics or clothes.\nThe present article of manufacture can further comprise a set of instructions to communicate methods of using the present compositions to the consumer. When articles of manufacture comprise devices or substrates that dispense the said composition, preferred devices or substrates will disperse the said composition in a uniform manner so as to minimize staining.\nThe present invention also comprises the use of small particle diameter droplets of the present compositions to treat fabrics, in order to provide superior performance, e.g., the method of applying the compositions to fabrics, etc. as very small particles (droplets) preferably having weight average diameter particle sizes (diameters) of from about 5 xcexcm to about 250 xcexcm, more preferably from about 10 xcexcm to about 120 xcexcm, and even more preferably from about 20 xcexcm to about 100 xcexcm."}
{"text": "Corrosion is generally a problem in any system in which ferrous metals are in contact with aqueous solutions. Corrosion is the electrochemical reaction of metal with its environment. Corrosion is a destructive reaction which, simply stated, is the reversion of refined metals to their natural state. For example, iron ore is iron oxide. Iron oxide is refined into steel. When the steel corrodes, it forms iron oxide.\nThe metallurgy of boiler systems are predominantly iron-containing steel in contact with high temperature (approximately 200.degree. F., under 35-350 psig of pressure) aqueous solutions. Most industrial boiler and feedwater systems are constructed of carbon steel and sometimes with copper alloy and/or stainless steel feedwater heaters and condensers. Some have stainless steel super-heater elements.\nCorrosion is one of the main causes of reduced reliability in boiler systems. The most common causes of corrosion are dissolved gases such as carbon dioxide and oxygen, under-deposit attack, low pH and attack of areas weakened by mechanical stress resulting in stress and fatigue cracking. Corrosion control in boilers varies with the type of corrosion encountered. Generally, corrosion in boilers may be controlled through procedures such as maintaining proper pH and alkalinity levels, reducing mechanical stresses, operating within design temperature and pressure specifications, by proper application of chemical corrosion inhibiting treatments and by controlling oxygen and boiler feedwater contamination. Generally boiler feedwater should be low in contaminants. The American Society of Mechanical Engineers (ASME) Consensus for Industrial Boilers specifies maximum levels of contaminants for corrosion and deposition control in boiler systems. The consensus is that feedwater oxygen, iron and copper content should be less than 7 parts per billion (ppb) oxygen, 20 ppb iron and 15 ppb copper for a 900 psig boiler and that pH should be maintained between 8.5 and 9.5 for boiler system corrosion protection.\nIt is commonly held that in boiler systems substantially free of oxygen (i.e., less than about 7 ppb oxygen) that ferrometal corrosion can be expressed in terms of the following equation: EQU Fe+2H.sub.2 O.fwdarw.Fe(OH).sub.2 +H.sub.2.\nTherefore, Fe.sup.+2 (ferrous ion) is a useful material to determine the corrosion rate of ferrometals in substantially oxygen free systems such as boilers. However, in the presence of oxygen the following reaction occurs: EQU 3Fe(OH.sub.2).fwdarw.Fe.sub.3 O.sub.4 +H.sub.2 +2H.sub.2 O\nMagnetite (Fe.sub.3 O.sub.4) is a complex of the hydroxides of iron (III) and iron (II) which can be expressed by the formula: EQU Fe.sup.3+ !.sup.IV Fe.sup.2+ Fe.sup.3+ !.sup.IV O.sub.4.\nOther ferric-ferrous oxide species also exist, however, under the conditions present in most boiler water, magnetite is the most common species. As used herein the term \"ferric-ferrous oxide\" means magnetite as well as other ferric-ferrous oxide species such as maghemite (.gamma.Fe.sub.2 O.sub.3) and feroxyhyte (.delta.FeOOH) which are ferrimagnetic. Ferric-ferrous oxides are often introduced into a boiler system by boiler feedwater. Ferric-ferrous oxides can also be produced within the boiler through reaction of ferrous ions with trace free oxygen and via electrochemical reactions.\nFerrion chelation chemicals are compounds which form complex undissociated cations with divalent metal ions. Such chelation chemicals include 1,10-phenanthrolinre; 4,7-diphenyl-1,10-phenanthroline; 2,2.sub.1 bipyridine; 2,6-bis(2-pyridyl)-pyridine; 2,4,6-(bis(2-pyridyl)-1,3,5-triazine; phenyl 2-pyridyl ketoxime; 3-(2-pyridyl)-5,6-bis(4-phenylsulfonic acid)-1,2,4-nitrilotriacetic acid and salts thereof. Certain of these ferrion chelation chemicals react with ferrous ions to form colored species (ferroin reagents) which are useful in the coloromertric determination of ferrous ion concentration. It is generally believed that ferroin reagents do not change color when the reagent complexes with ferric (Fe.sup.+3) ions in water. However, some research suggests that oxides, such as magnetite, may be captured by the ferroin reagents and thus bias ferrous ion content determinations (see Matijevic, E. Han-Chyen, Chang, Interactions of Metal Hydrous Oxides with Chelating Agents: Part V Ferric-ferrous oxides-EDTA, J. Colloid Interface Sci., Fin. Chem. Lett. (1982); and Bohnsak, G. The Solubility of Ferric-ferrous oxides in Water and in Aqueous Solutions of Acid and Alkali (Essen: Vielkan-Verlag, 1987)).\nTherefore, conventional techniques to measure ferrous ion content as an indicator of ferrous metal corrosion may inadvertently include some ferric-ferrous oxides concentration in the calculation of ferrous ion thereby yielding higher corrosion rate determinations then are actually occurring in the studied boiler system.\nThus a need exists for a method to remove ferric-ferrous oxides from a liquid containing ferric-ferrous oxides and ferrous ions.\nA need also exists for a method to measure the ferrous ion concentration in boiler water without contamination of the rate determination results due to the presence of ferric-ferrous oxides in the boiler water.\nIt is, therefore, an object of this invention to provide a method for removing ferric-ferrous oxides from a liquid containing ferric-ferrous oxides and ferrous ions.\nIt is also an object of this invention to provide an accurate method to determine the corrosion rate of ferrous metals exposed to a boiler water in a boiler system which contains ferrous ions, ferric ions and/or ferric-ferrous oxides."}
{"text": "1. Field of the Invention\nThe present invention relates to a radio communication system and a repeater, more particularly, to a radio communication system, which repeats a signal transmitted from a base station by a repeater that is multistage-connected by a leaky coaxial cable (LCX) and communicates with a mobile terminal.\n2. Related Art\nIn the field of a train radio communication and a PHS (Personal Handy-Phone System) base station installed in metropolitan railway stations, a radio communication system using a leaky coaxial cable as antenna has been put into practical use. In such a leaky coaxial cable, a slit for leaking a radio wave is provided at an outer conductor of a coaxial cable.\nThe radio communication system using the leaky coaxial cable is advantageous in that only one transmitter can cover a long distance communication area and a stable communication can be realized even in an environment with a lot of metallic reflecting objects or obstacles. Furthermore, the radio communication system using the leaky coaxial cable is advantageous in that the communication area can be limited to a peripheral area of the leaky coaxial cable, so that the radio wave cannot be leaked to an unnecessary area, thereby realizing a high security.\nFIG. 10 shows a conventional radio communication system using the leaky coaxial cable. Referring to FIG. 10, a signal line 104 from a base station 103 is connected to each pre-amplifier 102 to which a leaky coaxial cable 101 is connected, and a communication is carried out between the base station 103 and a mobile terminal 105 via the leaky coaxial cable 101.\nFIG. 11 shows another conventional radio communication system using the leaky coaxial cable, in which a plurality of repeaters 111 are multistage-connected by leaky coaxial cables 101, and each of the repeaters 111 amplifies signal, thereby increasing a communication area. Referring to FIG. 11, the graph shows a relationship between a transmission distance of signals and an electric power (i.e. intensity of the signals), in which L6 indicates an intensity of signals transmitted from the base station 103 and propagated through the leaky coaxial cable 101, L7 indicates an intensity of received signals of a mobile terminal 105 located at a point distant from the leaky coaxial cable 101 with a predetermined distance, and L8 indicates a minimum receiver sensitivity of the mobile terminal 105.\nIn the structure shown in FIG. 10, it is necessary to transmit the signal to be transmitted from the base station 103 to each pre-amplifier 102 by another cable (a signal line 104) which is different from the leaky coaxial cable 101, so that two kinds of cables i.e. the leaky coaxial cable 101 functioning as an antenna and the signal line 104 are required.\nOn the other hand, in the structure shown in FIG. 11, the leaky coaxial cable 101 functions as both an antenna and a signal line, so that the number of components can be reduced. Therefore, it is possible to reduce the cost compared with the structure shown in FIG. 10.\nJapanese Patent Laid-Open No. 2006-295433 (JP-A 2006-295433), Japanese Patent Laid-Open No. 11-186943 (JP-A 11-186943), and Japanese Patent Laid-Open No. 2004-328121 (JP-A 2004-328121) disclose examples of radio communication system using the leaky coaxial cable."}
{"text": "Conventional image detection techniques require scanning a picture in its entirety to detect a target image therein. However, the conventional image detection techniques necessitate a large number of scan windows and thus render the system busier to the detriment of its performance, not to mention that it is a waste of time detecting for non-target images. Referring to FIG. 1(A), a conventional image detection technique requires scanning a picture in its entirety to detect a target image therein and thus necessitates a large number of scan windows as shown in FIG. 1(B).\nIn an attempt to overcome the aforesaid drawback of the prior art, a speeded up robust features (SURF) algorithm was invented and involves retrieving characteristic points from a picture and then performing a block scan on the retrieved characteristic points to detect for a target image. SURF has an advantage: it speeds up image detection, because it does not require scanning the picture in its entirety. However, SURF has a drawback: the block scan performed on the retrieved characteristic points is disadvantaged by the presence of non-target image blocks. As a result, SURF still has room for improvement in terms of efficiency."}
{"text": "Since their introduction by Kass et al. (Kass 1988), a variety of deformable models have been proposed and used in a numerous areas such as image segmentation, shape modeling, and object tracking. Existing deformable models can be classified as either parametric deformable models (PDMs) (Kass, 1988; Cohen, 1991; Xu, 1998; McInery, 1996; Cottes, 1991, 2001; Shen, 2001; Liu, 2005) or geometric deformable models (GDMs) (Malladi, 1995; Caselles, 1997; Yezzi, 1997; Sethian, 1999; Han, 2003) according to their representation and implementation. In particular, PDMs are represented explicitly as parameterized curves or surfaces in a Lagrangian formulation. GDMS, on the other hand, are represented implicitly as level sets of higher-dimensional level set functions. That is, to model a contour within a 2-D space, or a surface within a 3-D space, we define a continuous function (a level set function) which takes a certain value (e.g. zero) on the contour or surface (referred to as the “level set”), and a non-zero value elsewhere. The level set function evolves according to an Eulerian formulation (Sethian 1999).\nGDMs are superior to PDMs in the following aspects:\nModel representation and construction. GDMs are completely intrinsic and independent of the parameterization of the evolving contour. Thus, there is no need to add or remove nodes from the initial parameterization or adjust the spacing of the nodes. For PDMs, particularly, in the 3D cases, it is not a trivial task to generate parameterized representations for surface models.\nModel property calculation. The intrinsic geometric properties of a geometric deformable model, such as its normal and curvature, can be easily determined from the level set function.\nShape adaptation. In GDMs, the propagating contour or surface can automatically change its topology (e.g., merge or split) without requiring an elaborate mechanism to handle such changes.\nSelf-intersection detection. When the topology preserving transformations are required for a deformable model, self-intersections can become a problem. In GDMs, digital topology methods (Bertrand 1994) are available to detect self-intersections, while in PDMs, the computational demands related to self-intersection detection are usually very high (MacDonald 2000).\nOn the contrary, PDMs have the following merits over GDMs:\nModel constraints. Since PDMs are represented globally, it is convenient to add high level constraints on shape and topology. These constraints are particularly important when local features of the model and image are liable to deform the model in unexpected ways.\nLeakage prevention. Leakage from the boundary is less common in PDMs, whereas it is sensitive to boundary gaps and weak edges in GDMs.\nConvergence detection. The deformation procedures in PDMs have an explicit stop condition which is the minimization of some kinds of energy function, whereas in GDMs, the procedure usually stops when the certain number of iterations reaches a predetermined maximum, but no simple and fast methods are available to automatically estimate a suitable maximum number of iterations."}
{"text": "1. Field of the Invention\nThe present application relates to an electronic apparatus provided with a slot to which an expansion module can be attached detachably.\n2. Description of Related Art\nRecent notebook computers are provided with a random access memory (RAM) that temporarily stores programs or data in order to execute various information processes at a central processing unit (CPU). As the size of a RAM increases, the size of an executable program or the number of simultaneously executable programs can be increased.\nTypically, a RAM is formed by mounting a memory chip on a board in the form of a module. When this memory module is attached to a memory slot provided in a notebook computer, the memory size of the notebook computer can be increased. This memory slot typically is sealed with a cover in order to prevent a foreign particle from entering from the outside. The cover is typically in the shape of a plate. Patent Documents 1 and 2 disclose plate-shaped covers. The covers disclosed in Patent Documents 1 (JP 2003-076439A) and 2 (JP H10-268976A) are attached to an apparatus main body through claw coupling or screw threading.\nHowever, in the case of a plate-shaped cover, when opening a memory slot in order to attach or detach a memory module, it is necessary to disconnect the claw coupling or screw threading, and then cause the cover to be detached by its own weight from an apparatus main body. More specifically, the electronic apparatus is initially positioned such that a face having the memory slot is oriented upward, claw coupling or screw threading is then disconnected, the electronic apparatus is then positioned such that the face having the memory slot is oriented downward, and the cover is detached by its own weight from the electronic apparatus. Subsequently, the memory module is attached or detached in a state where the face having the memory slot is oriented upward. Accordingly, when detaching a cover in order to attach or detach a memory module, it is necessary to change the posture of the electronic apparatus a plurality of times, which is very troublesome."}
{"text": "The present invention relates to programmable circuits. More particularly, the present invention relates to state machines in programmable logic devices.\nNon-volatile memory devices, such as EPROM, EEPROM, and Flash EEPROM, store data even after power is turned off. One common application of EEPROMs is in programmable logic devices (PLDs). PLDs are standard semiconductor components purchased by systems manufacturers in a xe2x80x9cblankxe2x80x9d state that can be custom configured into a virtually unlimited number of specific logic functions. PLDs provide system designers with the ability to quickly create custom logic functions to provide product differentiation without sacrificing rapid time to market. PLDs may be reprogrammable, meaning that the logic configuration can be modified after the initial programming.\nOne type of PLD is known as a Field-Programmable Gate Array (FPGA). An FPGA is a general purpose device that can be programmed by an end user to perform one or more selected functions. An FPGA typically includes an array of individually programmable logic cells (PLCs), each of which is programmably interconnected to other PLCs and to input/output (I/O) pins via a programmable routing structure to provide the selected function. Examples of such devices are exemplified in U.S. Pat. Nos. 4,642,487; 4,706,216; and 4,758,985.\nAn FPGA device can be characterized as an integrated circuit that may include four major features:\n(1) A user-accessible, configurable memory device, such as SRAM, EPROM, EEPROM, anti-fused, fused, or other, is provided in the FPGA device so as to be at least once-programmable by device users for defining user-provided configuration information. Static Random Access Memory or SRAM is a form of reprogrammable memory that may be differently programmed many times. Electrically Erasable programmable ROM or EEPROM is another example of nonvolatile reprogrammable memory. The configurable memory of an FPGA device may be formed of a mixture of different kinds of memory elements if desired (e.g., SRAM and EEPROM).\n(2) Input/Output Blocks (IOBs) are provided for interconnecting other internal circuit components of the FPGA device with external circuitry. The IOBs may have fixed configurations or they may be configurable in accordance with user-provided configuration information.\n(3) PLCs are provided for carrying out user-programmed logic functions (e.g., logic gates) as defined by user-provided configuration information. Typically, each of the many PLCs of an FPGA has at least one lookup table (LUT) that is user-configurable to define any desired truth table. A PLC may have other resources such as LUT input signal pre-processing resources and LUT output signal post-processing resources.\n(4) An interconnect network is provided for carrying signal traffic within the FPGA device between various PLCs and/or between various IOBs and/or between various IOBs and PLCs. At least part of the interconnect network is typically configurable so as to allow for programmably-defined routing of signals between various PLCs and/or IOBs in accordance with user-defined routing information.\nUsing these features, traditional FPGA architectures implement state machines by using the LUT-based logic elements in a decoding structure and as register storage elements. The state machine typically includes multiple states that are cycled through based on input signals. Typically, every clock cycle, the state of the input signals dictates whether the state machine maintains its current state or changes to a new state. Although state machines may be easily constructed in an FPGA, as the state machine grows large, performance begins to suffer as logic elements invariably get located at greater distances apart, thereby slowing propagation times. Additionally, such state machines consume a lot of area.\nThus, it is desirable to provide a PLD that can free PLC resources while increasing overall functionality. It is also desirable to implement an efficient state machine within a PLD.\nA PLD is disclosed that includes a state machine integrated into a block memory. The state machine includes state machine logic and memory elements from the block memory. The state machine logic and memory elements together may be used as an instruction unit of a processor. In such a case, the instruction unit functions as a counter, instruction store, and possibly a decode that is coupled to a processor execution unit to form a high-performance, embedded processor.\nIn one aspect, the state machine logic includes a program counter and an adder coupled to memory elements within the block memory.\nIn another aspect, control logic may be added to the state machine to receive state machine inputs and to determine the next state of the state machine.\nIn another aspect, the state machine may further include a stack pointer that is used to implement a stack within the block memory.\nIn yet another aspect, the state machine allows for implementing certain classes of instructions, such as jumps, branches and calls. Additionally, processor interrupts can be handled within the memory block.\nThese and other aspects will become apparent from the following detailed description, which makes reference to the accompanying drawings."}
{"text": "Turbines are conventionally started hydraulically, electrically or pneumatically. If hydraulic or pneumatic systems are used to start an aircraft auxiliary power unit (APU) electrical power is not normally available until the APU output shaft has exceeded some initial rotational speed to drive a generator. The consequences of not having initial electrical power include the inability to use electronic controller circuitry normally associated with APUs to perform such functions as opening and closing valves for the starting systems. Previously, when pneumatic or hydraulic starting systems were used, the operator had to manually close down the starting system after the APU began operation. This is not a reliable approach because if the starting system is turned off too early the APU will fail to complete the starting process; if, in the alternative, the starting system is turned off too late the starter motor may overspeed which makes failure of the starter system likely.\nTo address these and other problems, various previous methods have been devised. U.S. Pat. No. 3,812,378 to Coman describes the use of a permanent magnet generator for providing a first pulse train with each of the pulses having a duration that is a function of the rotational speed of the starter rotor. A second pulse generator is also provided, with the second pulse train having a constant pulse width; the second pulse train time width corresponds to a predetermined rotational speed of the starter rotor at which the starter is to be turned off. The first and second pulse trains are provided to a comparator circuit that produces an output signal when the duration of a pulse from the first pulse train is shorter than a pulse from the second pulse train, i.e. the starter rotational speed exceeds the predetermined rotational speed. This output signal from the comparator circuit is then used to turn the starter off. The clear disadvantage of this system is the requirement for electrical power from the very initiation of the starting process.\nAnother starter control system is described in, U.S. Pat. No. 4,747,270 to Klie et al, here a pneumatic starting device is disclosed. The pressure in a reservoir supplying gas to a pneumatic starting device is monitored, and the monitored pressure is exclusively used to end the starting process. As the gas pressure drops in the reservoir the maximum speed of the starter is limited. Since pressure drop in the reservoir is used to end the starting procedure, the reservoir must be constructed with only a capacity sufficient for carrying out one starting. This is an inescapable disadvantage, and is especially problematic for aircraft operating from austere airfields."}
{"text": "1. Field of the Invention\nThe present invention relates to a pedal retraction amount control apparatus for a running vehicle.\n2. Description of the Related Art\nAs a pedal retraction amount control apparatus for a running vehicle which prevents a pedal from moving in a rearward direction with respect to a vehicle body, and moving into a passenger compartment side at the time of a vehicle head-on collision, there is, for example, the apparatus disclosed in Japanese Patent Publication No. 3239790. In the apparatus disclosed in this publication, when at the time of a vehicle head-on collision, the vehicle body is deformed so that the whole pedal system moves in the rearward direction of the vehicle, a steering hanger which is provided on a horizontal beam member connecting between pillars restricts the movement of this pedal in the rearward direction beneath the support shaft of the pedal. As a result, the pedal rotates about the support shaft so that the tread portion moves towards the front side, and is thus moved from the passenger compartment in the forward direction.\nIn a pedal retraction amount control apparatus, the greater the amount of forward movement of the pedal the better. However, in the abovementioned case, the amount of forward movement of the pedal is mainly determined by the rotation amount of the pedal about the support shaft, and there are cases where this is not sufficient."}
{"text": "1. Field of the Invention\nThe present invention relates to a knee protecting airbag device in which an airbag to be expanded and inflated with an inflating gas let in protects the knee of a driver.\n2. Description of Related Art\nIn the prior art, the device for protecting the knee of a driver is known in Japanese Patent Laid-Open No. 6-32195, EP Publication No. 0684167A1, EP Publication No. 0818360A1 and so on. In the prior art, more specifically, there were a device in which a pad is arranged near a column cover below a steering wheel to protrude when activated to hold the knee of a driver, or a device having an airbag which is expanded and inflated when activated to protect the knee of the driver.\nIn the knee protecting devices of the prior art, however, the pad or airbag is protruded into a narrow space between the column cover and the knee of the driver. Especially at the action time, the driver is mostly depressing a brake pedal. In this case, the space between the column cover and the knee of the driver is further narrowed.\nTherefore, devices of the prior art had a room for improvement in proper protection of the knee of a driver."}
{"text": "1. Field of the Invention\nThe invention relates to a speaker assembly. More particularly, the invention relates to an articulating speaker assembly for integration within a passenger service unit of an aircraft.\n2. Description of the Related Art\nThe current global community has made it possible for people all around the country, and around the world, to interact for both business and personal reasons. For many people, this requires they spend considerable time traveling from one location to another location. More often than not, these people travel in aircraft. Whether these people travel in private or commercial aircraft, they desire high quality entertainment during the many hours they spend within the confines of an aircraft. However, while high quality entertainment, for example, digital video with CD quality sound, is readily available for theater and home use, the weight, size and available space requirements for use in aircraft make it very difficult to incorporate high fidelity systems within an aircraft. This problem is especially pronounced for audio loudspeaker assemblies when one attempts to meet the size, weight and shape requirements for use in aircrafts.\nIn the aircraft industry, great priority is placed upon component weight and size reduction. In addition, spacing and positioning of the loudspeaker assemblies are a great priority to those optimizing the operation of aircraft. The size, weight and shape of conventional terrestrial loudspeaker assembly designs adversely affect range and payload. These concerns are notable when one attempts to make changes within smaller, private jets. For example, a small increase in the weight carried by an aircraft results in a substantial increase in fuel consumption of the aircraft. In addition, the limited space available within an aircraft dictates the use of any space within the aircraft be carefully considered by those responsible for ensuring the comfort of passengers.\nLightweight and compact audio loudspeakers are currently available. These loudspeakers, however, substantially compromise sound quality for reductions in size and weight. An individual wishing to add an audio system to an aircraft must make a choice between high fidelity loudspeakers not suiting the size and weight requirements of the aircraft and lower quality loudspeakers providing desirable size and weight characteristics.\nA need, therefore, exists for a loudspeaker assembly providing high fidelity sound, while meeting the size, weight and positioning profile requirements of an aircraft. The present invention provides such a loudspeaker assembly."}
{"text": "The present invention relates to the arrangement of a plurality of solar cell elements on a supporting substrate.\nRecently, some kinds of electronic apparatus such as electronic wristwatches have been developed wherein solar cells are employed as power supply means.\nThe electronic energy generated from a solar cell device will increase when the size of a light receiving surface of the solar cell device is increased. Therefore, it is of great importance to increase the light receiving area of the solar cell device without increasing the total size of the solar cell device.\nAccordingly, an object of the present invention is to enhance the opto-electric converting efficiency of a solar cell device including a plurality of solar cell elements.\nAnother object of the present invention is to provide a novel arrangement of a plurality of solar cell elements for enhancing the opto-electric converting efficiency with a minimum package size.\nOther objects and further scope of applicability of the present invention will become apparent from the detailed description given hereinafter. It should be understood, however, that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description.\nTo achieve the above objects, pursuant to an embodiment of the present invention, a plurality of solar cell elements are arranged on a supporting substrate. Each solar cell element has a light receiving surface, where electrodes are formed which are connected to a P-type region and an N-type region contained within the solar cell element, respectively. The electrodes formed on the light receiving surfaces of the plurality of solar cell elements arranged on the supporting substrate are electrically connected to each other via light transmitting, electrically conductive wiring means formed on the light receiving surfaces, thereby developing electric energy."}
{"text": "The described embodiments relate generally to portable listening devices such as earbuds and other types of in-ear listening devices, and to cases for storing and charging such devices.\nPortable listening devices can be used with a wide variety of electronic devices such as portable media players, smart phones, tablet computers, laptop computers and stereo systems among others. Portable listening devices have historically included one or more small speakers configured to be place on, in, or near a user's ear, structural components that hold the speakers in place, and a cable that electrically connects the portable listening device to an audio source. Other portable listening devices can be wireless devices that do not include a cable and instead, wirelessly receive a stream of audio data from a wireless audio source.\nWhile wireless portable listening devices have many advantages over wired devices, they also have some potential drawbacks. For example, wireless portable listening devices, typically require a battery, such as a rechargeable battery, that provides power to the wireless communication circuitry and other components of the device. For many currently available wireless portable listening devices, charge can be restored to the rechargeable battery in the device by physically connecting the portable listening device to a power source, which typically requires that the wireless portable listening device have a pair of electrical contacts to receive the charge.\nIn many devices such electrical contacts are positioned within a receptacle connector in the wireless portable listening device. The receptacle connector typically includes a cavity in the wireless portable listening device that provides an avenue within which dust and moisture can intrude and damage the device. Furthermore, a user of the electronic device has to physically connect the charging cable to the receptacle in order to charge the battery. Some other wireless portable listening devices including charging contacts at an external surface of the device, which is a significant improvement with respect to moisture resistance and other potential problems as compared to internal contacts positioned within a receptacle connector cavity. Even external contacts, however, can still result in potential paths for moisture ingression under extreme conditions (for example, at the seams between the contacts and housing) and/or potential corrosion of the contacts when repeatedly exposed to corrosive liquids, such as sweat."}
{"text": "This invention relates generally to making of masks for semiconductor use and more particularly to calibration of an electron beam tool used in writing to a substrate.\nIn today\"\"s fabrication of Integrated Circuits (IC) and other semiconductor devices, lithographic delineation procedures are used to yield positive or negative images to bring about selective processing, e.g. etching, implantation, diffusion, deposition, etc. This is especially true in fabrications of masks where the fabrication tool provides Blocking regions and Transparent regions which when illuminated by electron radiation yields an image defined by relatively low and high electron intensities, respectively. A Blocking region is usually defined as the mask region resulting in a degree of electron attenuation in the image which is of consequence in device fabrication. By contrast, a Transparent region is the mask region resulting in a degree of electron attenuation in the image which is small relative to blocking regions in terms of device fabrication.\nIn the semiconductor industry, there is a continuing trend toward an increased device density. To achieve this, there is a continued effort towards the scaling down of device dimensions on semiconductor wafers. As smaller feature sizes become the new requirements (i.e. decreased width and spacing of interconnecting lines, etc.), new ways have to be utilized to achieve their manufacturing. High resolution lithographic processes are used as one of these manufacturing techniques to yield small component features.\nIn general, lithography refers to processes for pattern transfer between various media. In lithography for integrated circuit fabrication, a silicon slice, the wafer, is coated uniformly with a radiation-sensitive film, the resist. The film is then selectively exposed to radiation, such as optical light, x-rays, or an electron beam, through an intervening master template or the mask, forming a particular pattern. (In a mask, this leads to the creation of Blocking and Transparent regions which when illuminated by electron radiation yields an image defined by relatively low and high electron intensities, respectively.)\nMost often exposed areas of the coating become either more or less soluble than the unexposed areas (depending on the type of coating) in a particular solvent developer. The more soluble areas are removed with the developer in a developing step, the less soluble areas remain on the silicon wafer forming a patterned coating. The pattern corresponds to the image of the mask or its negative. The patterned resist is used in further processing of the silicon wafer.\nAt various stages in forming the patterned resist coating and processing the silicon wafer, it is desirable to measure critical dimensions resulting from the lithographic process. Critical dimensions can include the size of features in the wafer or patterned resist such as line widths, line spacing, and contact dimensions. Several calibration methods are developed and can be used such as scanning electron microscopy (SEM) systems.\nIn such calibration system, because of the super fine structures to be calibrated, an electron beam is often scanned across the sample. The beam interacts with the sample to produce measurable responses that vary with position over the course of a scan. There are several other ways to check for accuracy as well to ensure image precision at a later time.\nOne of the methods used historically by the E beam lithography tools is one that involves checking for stitch errors in determining more important underlying errors. This is often done by comparing the measurement of 4 corners of a target taken in a grid like subfield and comparing it with the measurement of the same corners taken from adjacent subfields. The values are usually then compared against a narrow and wide range for functional purposes. There are several problems with this approach, however. For one, the narrow range is easily exhaustible (exceeded), which makes continuous testing difficult. Second, this approach leaves much information uncovered which may mean that certain electromechanical system errors can go undetected. In order to make this approach more reliable, sometimes an operator has to be dispatched to visually and continuously monitor the testing which adds both inconvenience and cost to the functioning of the test. Therefore, an improved method and structure is needed to detect electromechanical problems in an automated manner and with more reliability\nThese and other objects are provided by the present invention for an apparatus and method for detection of electromechanical and mechanical errors in an electron beam device. First, a gridlike subfield is provided where each grid can be considered a target subfield. Then, a plurality of target points on each target subfield is indicated and the combined variance of displacement of target points on each target subfield is calculated to nominal position of the target subfield. Later, the combined target subfields stitching standard deviation or variance is also calculated by determining separation of each target point in each subfield with that of overlapping target points on the adjacent subfields. Variance statistics based on the sample variance of all measurements done in the field, variance statistics based on of measurements on a particular row and variance statistics based on differences between horizontally and vertically stitched measurements are made.\nA stitched measurement is a measurement made on a particular target two different ways: in this case, the same measurement made two times on two different raster scans. Ratios of these variances can be used to form significance tests related to analysis of variance (ANOVA) and in the form of F statistics, and are well known in statistical literature. A significance test is then conducted using data from the variances derived from adjacent sets of horizontal and vertical values of said target points resulting in calculation of separated vertical and horizontal FSTITCH values. When an F statistic, such as FSTITCH is greater than some threshold, it indicates statistical significance where significance tables for the F statistics are readily available in the literature for threshold selection depending on an acceptable overkill rate known in the literature as alpha significance. By comparing FSTITCH values for horizontal and vertical values to some threshold greater than one, an error alert can be provided when FSTITCH values exceed that threshold. The horizontal statistic and vertical statistics may be used to indicated isotropy or anisotropy. Lastly, a general F statistic based on average row variance and the variance over the entire field variance may be used as one indication of an electromechanical disturbance in raster scanned data, particularly in an E beam system."}
{"text": "1. Field of the Invention\nThe present invention relates to an ESD (electrostatic discharge) protection component and a method for manufacturing the EST protection component.\n2. Related Background Art\nThere is a known ESD protection component provided with an element body in which a plurality of insulator layers are stacked, a coil constructed by connecting a plurality of internal conductors to each other and arranged in the element body, and an ESD suppressor arranged in the element body and configured including first and second discharge electrodes arranged as separated from each other (e.g., cf. Japanese Patent Application Laid-Open Publication No. 2003-123936 (which will be referred to hereinafter as Patent Literature 1)). There is another known ESD protection component provided with an ESD suppressor configured including first and second discharge electrodes arranged as separated from each other, and a discharge inducing portion kept in contact with the first and second discharge electrodes so as to connect mutually opposed portions of the first and second discharge electrodes to each other and containing metal particles, in which a cavity portion is arranged so as to be in contact with the foregoing mutually opposed portions of the first and second discharge electrodes and with the discharge inducing portion (e.g., cf. Japanese Patent Application Laid-Open Publication No. 2011-243896 (which will be referred to hereinafter as Patent Literature 2))."}
{"text": "Phenyl ureas such as N'-3,4-dichlorophenyl-N,N-dimethyl-urea are known to possess herbicidal activity. However, the very broad spectrum herbicidal activity of these known phenyl ureas imposes severe limitations on the practical use of these compounds due to their lack of selectivity and potential injury to important crops such as corn and soybeans."}
{"text": "1. Field of the Invention\nThe present invention relates to a technology for controlling driving of a motor of a storage medium, depending on a rotational state of the motor.\n2. Description of the Related Art\nIn a magnetic disk device, a spindle motor is used to drive a magnetic disk. Two driving mode are used when driving the magnetic disk. One is a start-up mode in which the magnetic disk is driven from a non-rotating state and the other is a spin-up mode in which an already rotating magnetic disk is driven still faster.\nSpecifically, the start-up mode is a normal mode by which a spindle motor is driven. Through the start-up mode, a start-up control is performed for implementing control of phase switching (commutation) by which an electric current flows, from different directions, through coils forming the spindle motor, and when the spindle motor starts to rotate adequately, a spin-up control (or a BEMF(back electromotive force) control) is performed based on BEMF signals produced by the spindle motor.\nThe spin-up mode, on the other hand, includes driving the spindle motor after the magnetic disk is activated from a sleep state. In this case, the BEMF signals can be detected as the spindle motor is already in an inertial rotating state, and the BEMF control can be directly performed without going through the start-up control step.\nVarious technologies have been proposed to assuredly drive the spindle motor as well as to reduce a start-up time of the spindle motor. For example, Japanese Patent Application Laid-open No. H7-226016 discloses a technology for reducing the start-up time in the start-up mode, by stopping a steadily rotating motor, checking relative positions of a rotor and stator members, restarting the motor, and performing synchronous control. Specifically, the start-up time for the next Spin-On operation (at a time of starting up rotation) can be reduced by performing a brake operation and an aligning operation during Spin-Off operation (at a time of stopping rotation).\nHowever, when the spindle motor is not rotating adequately for proceeding an operation to the spin-up mode. (to perform a BEMA control), a longer start-up time is required compared to a time required when the spindle motor is rotating adequately, because it is necessary to retry the start-up operation repeatedly and perform a start-up control after the spindle motor is stopped.\nWith the conventional technology described above, although the start-up time is reduced for the Spin-On operation (at a time of starting up the rotation) only in the start-up mode, operation time cannot be reduced in the spin-up mode."}
{"text": "The present invention relates to a digital broadcast receiving and reproducing apparatus having a function of protecting the copyright by preventing unauthorized viewing or copying of digital broadcast.\nA digital broadcast receiving and reproducing apparatus according to a prior art is disclosed, for example, in Japanese Laid-open Patent No. H9-93561.\nFIG. 8 is a block diagram of digital broadcast receiving and reproducing apparatus of prior art.\nThis receiving and reproducing apparatus comprises a smart card (ID card) 804 storing a copy protection control code for judging permission of reception of broadcast, a digital VCR 805 recording digital data, and a digital broadcast receiver 801. The digital broadcast receiver 801 includes a descrambler 802 for canceling the scramble coded into program data, and a recording and digital output controller 803 for generating and issuing digital data and copy protection control code. The digital VCR 805 is\nconnected to the receiver 801 through data input and output ports\nand provided with control signals for input and output data and erasing data recorded in tape.\nIn the digital broadcast receiving and reproducing apparatus according to such prior art, as the memory device of digital data, it is expected to be developed not only into the digital VCR, but also into magnetic disk device, write once type magneto-optical disk device, and rewritable magneto-optical disk device. Moreover, there is an increasing demand for digital broadcast receiving and reproducing means for pay broadcast and/or pay-per-view, capable of protecting the copyright of the broadcast program by allowing only the subscribers to reproduce in ordinary pay broadcast or allowing to reproduce only once in pay-per-view, and without recording the copy protection control code, which is the key of broadcasting system, in the recording medium so as to be high in hacking resistance.\nTo solve the problems, the digital broadcast receiving and reproducing apparatus of the invention is a digital broadcast receiving and reproducing apparatus for receiving digital broadcast, demodulating received digital data, and taking out analog signal and copy protection information from the digital data, which comprises an apparatus individual key generator for generating and storing an individual key for the receiving and reproducing apparatus, a program key generator for generating a program key for every received program and storing this program key and the number of times of limited use in the copy protection information, a scrambler for multiplexing the program key and the demodulated digital data and scrambling by the apparatus individual key, a record judging unit for judging permission of recording of the scrambled digital data, and a digital data memory device for recording digital data through the record judging unit.\nThis digital broadcast receiving and reproducing apparatus further comprises a descrambler for descrambling the scrambled digital data reproduced from the memory device by using the apparatus individual key, extracting the program key from the scrambled digital data, and comparing the extracted program key and the program key stored in the program key generator. The digital broadcast receiving and reproducing apparatus of the invention is characterized by determining permission of output of digital data from the number of times of limited use and the result of the comparison.\nThus scrambled video and audio data can be reproduced, in the case of ordinary pay broadcast, only in the apparatus used for receiving or can be reproduced by the specified number of times in the apparatus used for receiving only, or in the case of pay-per-view, it is permitted to reproduce once only in the received apparatus. In this way, the key of the broadcast system, that is, the copy protection control code is not sent out of the receiving and reproducing apparatus, it has a high resistance to hacking, and not limited to the digital VCR, the digital broadcast receiving and reproducing apparatus is applicable to other recording device such as magnetic disk device, write once type magneto-optical disk device, or rewritable magneto-optical disk device."}
{"text": "The invention relates to management of talk groups in a telecommunications system at the dispatcher station of the telecommunications system.\nThe method according to the invention is intended for use particularly in PMR (Private Mobile Radio or Professional Mobile Radio) networks, i.e. in trunking networks, which are typically company networks or authority networks in which all channels are used by one or more companies or authorities. In these networks the subscribers have, in addition to their personal subscriber numbers, group numbers which indicate the group call group (talk group) to which the subscriber in question belongs so that calls intended for all subscribers of the group concerned can be transmitted to the subscribers of this group.\nGroup call is one of the most important functions of the PMR network. The group call is used e.g. in all kinds of activities in which several persons participate, particularly when the whole group needs to be aware of the course of events all the time. The group call is a conference call during which all participants may speak in turn and listen to each other. In group calls the whole group is called with one call number.\nHere a call or a group call refers to complete information exchange between two or more parties. The call may consist of one or more turns.\nOne example of a digital radio telephone or a mobile communication system in which the present invention can be applied is the TETRA system (Terrestrial Trunked Radio). In the TETRA system management of talk groups belongs to the responsibilities of a dispatcher. The dispatcher can perform different control functions related to talk groups and individual talk group users, i.e. subscribers. The dispatcher may, for example, create talk groups, change the composition of talk groups, add users to the system or change the users\"\" priority information. The dispatcher has a dispatcher workstation DWS at his disposal from which he receives all the information needed for managing talk groups. The workstation usually comprises a display and sound reproducers, e.g. headphones.\nThe problem associated with the system described above is that the dispatcher must look at the display to find out the group (or user) to which the talk heard at a given time belongs if the dispatcher does not recognize the speaker because the voices of all talk groups and subscribers are heard from the same source. One prior art solution to this problem is to use several sound reproducers, e.g. loudspeakers, so that the talk of one group is heard from a certain loudspeaker. The problem of this arrangement is that it requires several loudspeakers. Furthermore, if one wants to change the place from which the sound is heard, the physical location of the loudspeakers has to be changed.\nThe object of the invention is to provide a method and an arrangement implementing the method to solve the above-mentioned problems. The objects of the invention are achieved with a method of managing talk groups of a telecommunications system at the dispatcher station of the telecommunications system which comprises one or more talk groups which may consist of one or more users and which are controlled by a dispatcher at the dispatcher station, the system being characterized by creating an artificial acoustic space at the dispatcher station by means of a two-channel or a multichannel sound reproducing system and reproducing the voices of each talk group from a certain position of the acoustic space so that the dispatcher can recognize the talk group to which the voice belongs on the basis of the location of the voice. The preferred embodiments of the invention are disclosed in the dependent claims.\nThe invention is based on the idea of creating an acoustic impression of space artificially by means of a two-channel or a multichannel sound reproduction system, e.g. by loudspeakers or headphones, so that the voices belonging to different talk groups or users are heard from different positions at the dispatcher station. The acoustic space can be created e.g. as a three-dimensional space in which case voices can be arranged so that they are heard from any point of the dispatcher station.\nAn advantage of the method and arrangement of the invention is that the dispatcher can conclude the talk group or the user to which/whom the voice belongs on the basis of the virtual location of the voice. Thus the dispatcher does not need to look at the display to obtain this information.\nThe invention also relates to an arrangement for managing talk groups of a telecommunications system at the dispatcher station of the telecommunications system which comprises one or more talk groups which consist of one or more users and which are controlled by a dispatcher at the dispatcher station, the arrangement being characterized in that it comprises a two-channel or a multichannel sound reproducing system which is arranged to create an artificial acoustic space at the dispatcher station and reproduce the voices of each talk group so that they are heard from a certain point of the acoustic space and the dispatcher can recognize the talk group to which the voice belongs on the basis of the location of the voice. This kind of arrangement provides a simple solution for achieving the advantages of the method of the invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a vitreous silica crucible and a method of manufacturing a silicon ingot.\n2. Description of Related Art\nIn general, silicon single crystal is manufactured by melting high-purity polycrystalline silicon in a vitreous silica crucible to obtain silicon melt, dipping an end of a seed crystal to the silicon melt, and pulling the seed crystal while rotating it.\nA vitreous silica crucible is a crucible made of vitreous silica, and, in general, has two-layer structure of an inner layer of a synthetic vitreous silica layer (hereinafter, referred to as “synthetic layer”), and an outer layer of a natural vitreous silica layer (hereinafter, referred to as “natural layer”). The synthetic layer is a vitreous silica layer formed by fusing and solidifying chemically synthesized amorphous or crystalline silica (silicon oxide) powder, and the natural layer is a vitreous silica layer formed by fusing and solidifying silica powder obtained from natural mineral whose main component is a-quartz. The synthetic layer has very high purity, and thus it is possible to reduce the amount of impurities mixed in the silicon melt by providing the synthetic layer as the inner layer. The natural layer has high strength, and thus it is possible to enhance strength of the crucible by providing the natural layer as the outer layer.\nThe melting point of silicon is 1410 deg. C., and thus the temperature of silicon melt is kept at a temperature higher than 1410 deg. C. At such temperature, a vitreous silica crucible reacts with silicon melt, and the thickness of the crucible wall gradually decreases. When the thickness of the crucible wall decreases, the strength of the crucible is lowered. This leads to problems such as buckling and sidewall lowering of the crucible.\nIn order to solve such problems, there is known a technique to provide a layer, between the natural layer and the synthetic layer, for promoting crystallization of vitreous silica (e,g, JP Patent 3798907). When such a layer is provided and the crucible is heated for a long time, crystallization of the synthetic layer, which is in contact with silicon melt, is promoted, and eventually the entire portion of the inner layer of the crucible is crystallized. Crystalline silica has higher strength per unit thickness than vitreous silica, and, in addition, has lower reactivity with silicon melt than vitreous silica. Therefore, the crystallization enhances strength per unit thickness and suppresses reduction of the wall thickness of the crucible.\nFurthermore, there is known a technique to provide a layer, on the outside of the crucible, to promote crystallization. When such a layer is provided, the outer layer of the crucible is crystallized and the crucible strength is improved (e.g. JP-A-2000-247778).\nBy use of any of these techniques, the crucible strength is improved, and thus problems such as buckling and sidewall lowering of the crucible can be solved to some extent."}
{"text": "1. Field of the Invention\nThe present invention relates, in general, to the preparation of amino carboxylic acids, salts, and esters, and, in a preferred embodiment, to the preparation of N-(phosphonomethyl)glycine, its salts, and its esters, wherein the method of preparation comprises a carboxymethylation step.\n2. Description of Related Art\nAmino carboxylic acids are useful in various applications. Glycine, for example, is widely used as an additive in processed meat, beverages, and in other processed food stuffs. It is also used widely as a raw material for pharmaceuticals, agricultural chemicals, and pesticides. N-(phosphonomethyl)glycine, also known by its common name glyphosate, is a highly effective and commercially important herbicide useful for combating the presence of a wide variety of unwanted vegetation, including agricultural weeds. Between 1988 and 1991, approximately 13 to 20 million acres per year worldwide were treated with glyphosate, making it one of the most important herbicides in the world. Convenient and economical methods of preparing glyphosate and other amino carboxylic acids are, therefore, of great importance.\nFranz, et al. in Glyphosate: A Unique Global Herbicide (ACS Monograph 189, 1997) at p. 233-257 identify a number of routes by which glyphosate can be prepared. According to one of these, iminodiacetic acid disodium salt (DSIDA) is treated with formaldehyde and phosphorous acid or phosphorous trichloride to produce N-(phosphonomethyl)-iminodiacetic acid and sodium chloride. A carboxymethyl group on the N-(phosphonomethyl) iminodiacetic acid is then oxidatively cleaved in the presence of a carbon catalyst to produce glyphosate acid. A significant drawback of this method is that it produces as a side product three equivalents of sodium chloride per equivalent of glyphosate. Sodium chloride streams of this nature are difficult to recycle because typically after precipitation the salt contains significant quantities of entrapped organic matter. Such entrapped organic matter prevents the sodium chloride from being used for many purposes, for example in foods or feed. Further recrystallization of the sodium chloride adds cost which makes recycle economically unfeasible. Alternate methods of disposing of sodium chloride without detriment to the environment are expensive and difficult.\nFranz et al. (at 242-243) describe another method in which N-isopropylglycine is phosphonomethylated to produce N-isopropyl-N-(phosphonomethyl)glycine. In this method, the N-isopropyl-N-(phosphonomethyl)glycine is heated to 300xc2x0 C. with 50% sodium hydroxide and then treated with hydrochloric acid to produce glyphosate. The severe and costly conditions necessary to cleave the N-isopropyl group represents a significant disadvantage of that method. In addition, this method also produces a significant sodium chloride waste stream.\nIn U.S. Pat. No. 4,400,330, Wong discloses a method for the preparation of glyphosate in which 2,5-diketopiperazine is reacted with paraformaldehyde and a phosphorous trihalide in a carboxylic acid solvent to produce N,Nxe2x80x2-di(phosphonomethyl)-2,5-diketopiperazine. The product is then saponified to form a glyphosate sodium salt. The Wong method is limited by the fact diketopiperazine is a relatively expensive starting material. Furthermore, the conversion of glyphosate sodium salt to the acid form or to other salts produces an undesired sodium chloride waste stream.\nAmong the objects of the present invention, therefore, is the provision of a well-defined, low-cost process for the production of amino carboxylic acids, in general, and N-(phosphonomethyl)glycine, in particular, and the provision of such a process in which sodium chloride is not generated as a by-product.\nIn the process of the present invention, an N-acyl amino carboxylic acid is formed via a carboxymethylation reaction. In this reaction, a reaction mixture is formed which contains a base pair, carbon monoxide and an aldehyde with the base pair being derived from a carbamoyl compound and a carboxymethylation catalyst precursor. In a preferred embodiment, the carbamoyl compound and aldehyde are selected to yield an N-acyl amino carboxylic acid which is readily converted to N-(phosphonomethyl)glycine, or a salt or ester thereof having the following structure: \nwherein R7, R8, and R9 independently are hydrogen, hydrocarbyl, substituted hydrocarbyl, or an agronomically acceptable cation. In general, carbamoyl compounds which are selected to produce N-(phosphonomethyl)glycine correspond to structure (II): \nwherein R1 is hydrogen, hydrocarbyl, substituted hydrocarbyl, xe2x80x94NR3R4, xe2x80x94OR5, or xe2x80x94SR6;\nR2 and R2a are independently hydrogen, hydrocarbyl, or substituted hydrocarbyl;\nR3 and R4 are independently hydrogen, hydrocarbyl, or substituted hydrocarbyl; and\nR5 and R6 are independently hydrogen, hydrocarbyl, substituted hydrocarbyl, or a salt-forming cation;\nprovided, however, (1) at least one of R2 and R2a is hydrogen, hydroxymethyl, amidomethyl, or another substituent which, under the carboxymethylation reaction conditions, is capable of producing an Nxe2x80x94H bond, or (2) R1 is xe2x80x94NR3R4 and at least one of R3 and R4 is hydrogen, hydroxymethyl, amidomethyl, or another substituent which, under the carboxymethylation reaction conditions, is capable of producing an Nxe2x80x94H bond.\nIn one embodiment of the process of the present invention, therefore, an amino carboxylic acid or a salt or an ester thereof is prepared by carboxymethylation of a carbamoyl compound. In this process, a reaction mixture is formed by combining the carbamoyl compound and a carboxymethylation catalyst precursor in the presence of carbon monoxide and hydrogen. Water and an aldehyde are introduced into the reaction mixture after the carbamoyl compound and the carboxymethylation catalyst precursor are combined and the components of the reaction mixture are reacted to generate a product mixture containing an N-acyl amino carboxylic acid reaction product and a catalyst reaction product.\nIn another embodiment of the process of the present invention, a reaction mixture containing the carbamoyl compound, carbon monoxide, hydrogen, an aldehyde, and a carboxymethylation catalyst precursor derived from cobalt is formed. The components of the reaction mixture are reacted to generate a product mixture containing an N-acyl amino carboxylic acid reaction product and a catalyst reaction product. The catalyst reaction product is recovered from the product mixture and the catalyst reaction product is regenerated in the presence of the carbamoyl compound.\nIn a further embodiment, the process of the present invention is directed to the preparation of N-(phosphonomethyl)glycine or a salt or ester thereof. In this process, an N-acyl amino acid reaction product is prepared by carboxymethylating a carbamoyl compound in a reaction mixture formed by combining the carbamoyl compound, formaldehyde, carbon monoxide, hydrogen and a carboxymethylation catalyst precursor derived from cobalt. The N-acyl amino acid reaction product is converted to N-(phosphonomethyl)glycine or a salt or ester thereof wherein said conversion comprises deacylating the N-acyl amino acid reaction product to generate a carboxylic acid and an amino acid. The carboxylic acid is reacted with an amine to generate the carbamoyl compound or a compound from which the carbamoyl compound may be derived.\nIn a further embodiment, N-(phosphonomethyl)glycine or a salt or ester thereof is derived from N-acetyliminodiacetic acid. The N-acetyliminodiacetic acid is prepared by carboxymethylating acetamide in a reaction mixture formed by combining acetamide, acetic acid, water, formaldehyde, carbon monoxide, hydrogen, and a carboxymethylation catalyst precursor derived from cobalt. The N-acetyliminodiacetic acid is converted to N-(phosphonomethyl)glycine or a salt or ester thereof wherein said conversion comprises deacylating N-acetyliminodiacetic acid.\nIn a further embodiment, N-(phosphonomethyl)glycine or a salt or ester thereof is derived from an N-acyl amino acid carboxylic acid reaction product which is prepared from a reaction mixture containing a carbamoyl compound selected from among ureas and N-alkyl substituted amides, a carboxymethylation catalyst precursor, formaldehyde, and carbon monoxide. The N-acyl amino carboxylic acid reaction product is then converted to N-(phosphonomethyl)glycine or a salt or ester thereof. If the carbamoyl compound is an N-alkyl substituted amide, the conversion step(s) comprise oxidatively dealkylating the N-acyl amino carboxylic acid reaction product in the presence of oxygen using a noble metal catalyst.\nThe present invention is additionally directed to the certain key starting materials used and intermediates prepared in the process of the present invention.\nFor example, in one embodiment, the present invention is directed to a compound having the formula: \nIn another embodiment, the present invention is directed to a compound having the formula: \nIn yet another embodiment, the present invention is directed to an acetamide equivalent compound selected from the group consisting of compounds having the formula: \nwherein R13 and R14 are independently hydrogen, hydroxymethyl, alkyl, carboxymethyl, phosphonomethyl, or an ester or salt of carboxymethyl or phosphonomethyl; R15, R16 and R17 are independently alkyl or xe2x80x94NR3R4; and R3 and R4 are independently hydrogen, hydrocarbyl, or substituted hydrocarbyl.\nIn still a further embodiment, the present invention is directed to a compound having the formula: \nwherein R1 is hydrogen, hydrocarbyl, substituted hydrocarbyl, xe2x80x94NR3R4, or SR6; R3 and R4 are independently hydrogen, hydrocarbyl, or substituted hydrocarbyl; and R6 is hydrogen, hydrocarbyl, substituted hydrocarbyl, or a salt-forming cation.\nFurther scope of the applicability of the present invention will become apparent from the detailed description provided below. It should be understood, however, that the following detailed description and examples, while indicating preferred embodiments of the invention, are given by way of illustration only since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description."}
{"text": "As disclosed in British Pat. Nos. 1,539,733 (Sheldon et al.) and 2,055,799 (Thiault et al.), it is known that 3-phenoxybenzaldehydes are useful as pesticide intermediates and that they can be prepared from 3-bromobenzaldehyde acetals. Sheldon et al., who react their acetals with phenols or alkali metal phenolates in the presence of a copper catalyst, require an aprotic organic solvent for their reaction. Tiault et al. avoid the need for the aprotic solvent by reacting their acetal with an excess of phenol in the presence of a copper catalyst and potassium carbonate at a temperature of 170.degree.-210.degree. C.\nAlthough the process of Thiault et al. is advantageous in that it avoids several problems associated with the use of the solvents of Sheldon et al., it also has its disadvantages. The elevated temperatures required for the reaction lead to degradation of the acetal, and the potassium carbonate appears to poison the catalyst. Thus, the obtainable yield of product is comparatively low."}
{"text": "The present invention relates to cosmetic compositions containing, in combination, at least one polyorganosiloxane and an acrylic terpolymer, as well as to the use of these compositions for treating keratinous material.\nSilicones are cosmetic products which are particularly sought for their conditioning properties and their softening and disentangling properties, in particular in hair formulations.\nWhen certain silicones are used, such as those which are insoluble in aqueous medium, it is often necessary to introduce them in a thickened or gelled support in order to keep them in suspension and to obtain a stable composition.\nIt is also important to have available non-pasty, non-greasy compositions which spread well on the skin and the hair.\nWith the aim of obtaining a thickened or gelled support, thickening and/or gelling polymers are used. Gelling polymers are known which contain in their chain a hydrophilic part and a hydrophobic part consisting of a fatty chain, such as the product xe2x80x9cPemulen TR1xe2x80x9d sold by the company Goodrich or the xe2x80x9cAcrysolxe2x80x9d polymers sold by the company Rohm and Haas. The polymer xe2x80x9cPemulen TR1xe2x80x9d gives formulations of pasty texture which do not spread well. The polymer xe2x80x9cAcrysolxe2x80x9d does not have good gelling power and gives cloudy, unstable formulations.\nThe Applicant has discovered, surprisingly, that by using a novel family of thickening and/or gelling polymers and by combining them with polyorganosiloxanes, cosmetic formulations which have a satisfactory viscosity at a relatively low pH, which are stable, non-greasy, non-pasty and which spread well on the skin and the hair can be obtained.\nThe subject of the present invention is thus cosmetic compositions containing, in a cosmetically acceptable aqueous support, at least one polyorganosiloxane and an acrylic terpolymer which will be defined in greater detail later in the description.\nThis polymer makes it possible in particular to prepare rinse-out or leave-in, aqueous or aqueous-organic compositions containing cosmetically acceptable solvents, ranging from slightly gelled products to solid sticks or tubes.\nThe advantages of this terpolymer are that it is stable in electrolytic medium and has very good thickening power at a pH equal to or above 5.5, which allows a good level of viscosity to be achieved and allows high concentrations of alcohol to be used.\nWhen used in combination with at least one polyorganosiloxane, this polymer makes it possible to prepare gelled, non-pasty products which are easy to spread, feel soft when applied and are stable on storage.\nIt also makes it possible to improve the conditioning effect of silicones on the hair, in particular their feel, their softness and their disentangling ability.\nThe acrylic terpolymer used in accordance with the invention is soluble or swellable in alkalis. It is characterized in that it comprises:\na) about 20 to 70% by weight, preferably 25 to 55% by weight, of a carboxylic acid containing xcex1,xcex2-monoethylenic unsaturation;\nb) about 20 to 80% by weight, preferably 30 to 65% by weight, of a non-surfactant monomer containing monoethylenic unsaturation, which is different from a), and\nc) about 0.5 to 60% by weight, preferably 10 to 50% by weight, of a nonionic urethane monomer which is the product of reaction of a monohydric nonionic surfactant with a monoisocyanate containing monoethylenic unsaturation.\nThe carboxylic acid containing xcex1,xcex2-mono-ethylenic unsaturation a) can be chosen from many acids and in particular acrylic acid, methacrylic acid, itaconic acid and maleic acid. Methacrylic acid is preferred. A large proportion of acid is essential in order to give a polymer structure which dissolves and gives a thickening effect by reaction with an alkaline compound such as sodium hydroxide, alkanolamines, aminomethylpropanol or aminomethylpropanediol.\nThe terpolymer should also contain a large proportion, indicated above, of a monomer b) containing monoethylenic unsaturation which has no surfactant properties. The preferred monomers are those which give polymers that are water-insoluble when they are homopolymerized and are illustrated by C1-C4 alkyl acrylates such as methyl acrylate, ethyl acrylate and butyl acrylate, or corresponding methacrylates. The monomers more particularly preferred are methyl and ethyl acrylates. Other monomers which can be used are styrene, vinyltoluene, vinyl acetate, acrylonitrile and vinylidene chloride. Non-reactive monomers are preferred, such monomers being those in which the single ethylenic group is the only group which is reactive under the polymerization conditions. However, monomers which contain groups that are reactive under the action of heat can be used in certain situations, such as hydroxyethyl acrylate.\nThe monohydric nonionic surfactants used to obtain the nonionic urethane monomer c) are well known and are generally alkoxylated hydrophobic compounds containing an alkylene oxide forming the hydrophilic part of the molecule. The hydrophobes generally consist of an aliphatic alcohol or an alkylphenol in which a carbon chain containing at least six carbon atoms constitutes the hydrophobic part of the surfactant.\nThe preferred monohydric nonionic surfactants have the formula: \nin which R is a C6-C30 alkyl or C8-C30 aralkyl group, Rxe2x80x2 is a C1-C4 alkyl group, n is an average number ranging approximately from 5 to 150 and m is an average number ranging approximately from 0 to 50, with the condition that n is at least as large as m and that n+m=5-150.\nAs preferred C6-C30 alkyl groups, mention may be made of dodecyl and C18-C26 alkyl radicals. As aralkyl groups, mention may be made more particularly of (C8-C13)alkylphenyl groups. The preferred group Rxe2x80x2 is the methyl group.\nThe monoisocyanate containing monoethylenic unsaturation which is used to form the nonionic urethane monomer c) can be chosen from a wide variety of compounds. A compound containing any copolymerizable unsaturation such as acrylic or methacrylic unsaturation can be used. An allylic unsaturation imparted by allyl alcohol can also be used. The preferred monoethylenic monoisocyanate is xcex1,xcex1-dimethyl-m-isopropenyl-benzylisocyanate.\nThe acrylic terpolymer defined above is obtained by aqueous emulsion copolymerization of the components a), b) and c) which is entirely common and described in patent application EP-A-0,173,109.\nAs terpolymers which can be used according to the invention, mention may be made of the products of reaction of methacrylic acid as component a), of ethyl acrylate as component b) and of a nonionic urethane macromonomer as component c), having the following structure: \nin which pxe2x80x2 ranges from 6 to 150 and is preferably equal to 30 and R2 is a C8-C13 alkyl radical, such as that described in Example 3 of patent application EP-A-0,173,109.\nThe preferred acrylic terpolymer used according to the invention is obtained from methacrylic acid as component a), methyl acrylate as component b) and a nonionic urethane macromonomer as component c), having the following structure: \nin which p ranges from 6 to 150 and R1 is a C18-C26, alkyl radical preferably C20-C24, linear of plant origin, such as the docosyl radical.\nThe acrylic terpolymer is present in the cosmetic compositions of the invention in concentrations ranging from 0.01 to 20% by weight relative to the total weight of the composition, and preferably from 0.1 to 10% by weight.\nThe modified or non-modified polyorganosiloxanes used in the compositions according to the present invention are polyorganosiloxane oils or polyorganosiloxane gums or resins, in their natural state or in the form of solutions in organic solvents, or alternatively in the form of emulsions or microemulsions.\nAmong the polyorganosiloxanes used in accordance with the present invention, mention may be made, in a non-limiting manner, of:\nI. Volatile Silicones\nThese have a boiling point of between 60xc2x0 C. and 260xc2x0 C. They are chosen from cyclic silicones containing from 3 to 7 and preferably 4 to 5 silicon atoms. They are, for example, octamethylcyclotetrasiloxane sold under the name xe2x80x9cVolatile Silicone 7207xe2x80x9d by Union Carbide or xe2x80x9cSilbione 70045 V2xe2x80x9d by Rhxc3x4ne-Poulenc, decamethylcyclopentasiloxane sold under the name xe2x80x9cVolatile Silicone 7158xe2x80x9d by Union Carbide, xe2x80x9cSilbione 70045 V5 by Rhxc3x4ne-Poulenc, and mixtures thereof.\nMention is also made of cyclocopolymers of the dimethylsiloxane/methylalkylsiloxane type, such as xe2x80x9cVolatile Silicone FZ 3109xe2x80x9d sold by the company Union Carbide, which is a dimethylsiloxane/methyloctylsiloxane cyclocopolymer.\nII. Non-volatile Silicones\nThese consist mainly of:\n(i) polyalkylsiloxanes;\n(ii) polyarylsiloxanes;\n(iii) polyalkylarylsiloxanes;\n(iv) silicone gums;\n(v) silicone resins;\n(vi) organomodified polyorganosiloxanes;\n(vii) block copolymers having a polysiloxanepolyalkylene linear block as repeating unit;\n(viii) grafted silicone polymers containing a non-silicone organic skeleton, consisting of an organic main chain formed from organic monomers not containing silicone, onto which is grafted, within the said chain as well as, optionally, on at least one of its ends, at least one polysiloxane macromonomer;\n(ix) grafted silicone polymers containing a polysiloxane skeleton, grafted with non-silicone organic monomers, comprising a polysiloxane main chain onto which is grafted, within the said chain as well as, optionally, on at least one of its ends, at least one organic macromonomer containing no silicone;\n(x) or mixtures thereof.\nAmong the polyalkylsiloxanes, mention may be made mainly of linear polydimethylsiloxanes containing trimethylsilyl end groups, such as, for example, and in a non-limiting manner, the xe2x80x9cSilbionexe2x80x9d oils of the 70047 series sold by Rhxc3x4ne-Poulenc; the oil xe2x80x9c47 V 500 000xe2x80x9d from Rhxc3x4ne-Poulenc or certain xe2x80x9cViscasilxe2x80x9d oils from General Electric or xe2x80x9cMirasilxe2x80x9d oils from Rhxc3x4ne-Poulenc, and linear polydimethylsiloxanes containing hydroxydimethylsilyl end groups, such as the oils of the 48 V series from Rhxc3x4ne-Poulenc.\nIn this class of polyalkylsiloxanes, mention may also be made of the polyalkylsiloxanes sold by the company Goldschmidt under the names xe2x80x9cAbilwax 9800xe2x80x9d and xe2x80x9cAbilwax 9801xe2x80x9d, which are poly(C1-C20)alkylsiloxanes.\nAmong the polyalkylarylsiloxanes, mention may be made of linear and branched polymethylphenylsiloxanes, polydimethylmethylphenylsiloxanes and polydimethyldiphenylsiloxanes, such as, for example:\nthe oil xe2x80x9cRhodorsil 763xe2x80x9d from Rhxc3x4ne-Poulenc,\nthe oils xe2x80x9cSilbione 70641 V 30 and 70641 V 200xe2x80x9d from Rhxc3x4ne-Poulenc,\nthe product xe2x80x9cDC 556 Cosmetic Grade Fluidxe2x80x9d from Dow Corning,\nthe silicones of the PK series from Bayer, such as xe2x80x9cPK20xe2x80x9d,\nthe silicones of the PN and PH series from Bayer, such as xe2x80x9cPN 1000xe2x80x9d and PH 1000xe2x80x9d,\ncertain oils of the SF series from General Electric, such as SF 1250, SF 1265, SF 1154 and SF 1023.\nThe silicone gums in accordance with the present invention are polydiorganosiloxanes with a molecular mass of between 200,000 and 1,000,000, used alone or as a mixture in a solvent chosen from volatile silicones, polydimethylsiloxane (PDMS) oils, polyphenylmethylsiloxane (PPMS) oils, isoparaffins, methylene chloride, pentane, dodecane, tridecane and tetradecane, or mixtures thereof.\nMention is made, for example, of the following compounds:\npolydimethylsiloxane,\npoly[(dimethylsiloxane)/(methylvinylsiloxane)],\npoly[(dimethylsiloxane)/(diphenylsiloxane)],\npoly[(dimethylsiloxane)/(phenylmethylsiloxane)],\npoly[(dimethylsiloxane)/(diphenylsiloxane)/(methylvinylsiloxane)].\nMention may be made, for example, in a non-limiting manner, of the following mixtures:\n1) mixtures formed from a polydimethylsiloxane hydroxylated at the end of the chain (Dimethiconol according to the CTFA nomenclature) and from a cyclic polydimethylsiloxane (Cyclomethicone according to the CTFA nomenclature), such as the product xe2x80x9cQ2 1401xe2x80x9d sold by the company Dow Corning;\n2) mixtures formed from a polydimethylsiloxane gum with a cyclic silicone, such as the product xe2x80x9cSF 1214 Silicone Fluidxe2x80x9d from General Electric, which is an SE 30 gum of MW 500,000 dissolved in xe2x80x9cSF 1202 Silicone Fluidxe2x80x9d (decamethylcyclopentasiloxane);\n3) mixtures of two PDMSs of different viscosity, in particular a PDMS gum and a PDMS oil, such as the products xe2x80x9cSF 1236xe2x80x9d and xe2x80x9cCF 1241xe2x80x9d from the company General Electric. The product xe2x80x9cSF 1236xe2x80x9d is a mixture of an SE 30 gum defined above, with a viscosity of 20 m2/s, and of an SF 96 oil with a viscosity of 5xc3x9710xe2x88x925 m2/s (15% SE 30 gum and 85% SF 96 oil). The product xe2x80x9cCF 1241xe2x80x9d is a mixture of an SE 30 gum (33%) and of a PDMS (67%) with a viscosity of 10xe2x88x923 m2/s.\nThe polyorganopolysiloxane resins which can be used in accordance with the invention are preferably crosslinked siloxane systems containing the units: R2SiO2/2, RSiO3/2 and SiO4/2 in which R represents a hydrocarbon-based group having 1 to 6 carbon atoms or a phenyl group. Among these products, the ones which are particularly preferred are those in which R denotes a lower alkyl radical or a phenyl radical.\nAmong these resins, mention may be made of the product sold under the name xe2x80x9cDow Corning 593xe2x80x9d or those sold under the names xe2x80x9cSilicone Fluid SS 4230xe2x80x9d and\nAmong these resins, mention may be made of the product sold under the name xe2x80x9cDow Corning 593xe2x80x9d or those sold under the names xe2x80x9cSilicone Fluid SS 4230xe2x80x9d and xe2x80x9cSS 4267xe2x80x9d by the company General Electric and which are dimethyl/trimethylpolysiloxanes.\nThe organomodified silicones in accordance with the present invention are silicones as defined above, containing in their general structure one or more organofunctional groups directly attached to the siloxane chain or attached via a hydrocarbon-based radical.\nMention is made, for example, of silicones containing:\na) polyethyleneoxy and/or polypropyleneoxy groups, optionally containing alkyl groups, such as:\nthe product known as dimethiconecopolyol sold by the company Dow Corning under the name xe2x80x9cDC 1248xe2x80x9d, and alkyl (C12) methiconecopolyol sold by the company Dow Corning under the name xe2x80x9cQ2 5200xe2x80x9d,\nthe oils xe2x80x9cSilwetxe2x80x9d L 722, L 7500, L 77 and L 711 from the company Union Carbide,\nthe mixture of dimethiconecopolyol and of cyclomethicone, such as the product sold under the name xe2x80x9cQ2-3225Cxe2x80x9d by the company Dow Corning;\nthe product xe2x80x9cMirasil DMCOxe2x80x9d sold by Rhxc3x4ne-Poulenc;\nb) (per)fluoro groups, for instance trifluoroalkyl groups, such as, for example, those sold by the company General Electric under the names xe2x80x9cFF.150 Fluorosilicone Fluidxe2x80x9d or by the company Shin Etsu under the names xe2x80x9cX-22-819xe2x80x9d, xe2x80x9cX-22-820xe2x80x9d, xe2x80x9cX-22-821xe2x80x9d, xe2x80x9cX-22-822xe2x80x9d or xe2x80x9cFL 100xe2x80x9d;\nc) hydroxyacylamino groups, such as those described in European patent application EP-A-0,342,834, and in particular the silicone sold by the company Dow Corning under the name xe2x80x9cQ2-8413xe2x80x9d;\nd) thiol groups, such as in the silicones xe2x80x9cX 2-8360xe2x80x9d from Dow Corning or xe2x80x9cGP 72Axe2x80x9d and xe2x80x9cGP 71xe2x80x9d from Genesee; Union Carbide or the silicone known as xe2x80x9cAmodimethiconexe2x80x9d in the CTFA dictionary;\nf) carboxylate groups, such as the products described in European patent EP 186,507 from Chisso Corporation;\ng) hydroxylated groups, such as the polyorganosiloxanes containing a hydroxyalkyl function, described in patent application FR-A-2,589,476, and in particular polyorganosiloxanes containing a xcex3-hydroxy-propyl function;\nh) alkoxylated groups containing at least 12 carbon atoms, such as the product xe2x80x9cSilicone Copolymer F 7551xe2x80x9d from SWS Silicones and the products xe2x80x9cAbilwax 2428xe2x80x9d, xe2x80x9cAbilwax 2434xe2x80x9d and xe2x80x9cAbilwax 2440xe2x80x9d from the company Goldschmidt;\ni) acyloxyalkyl groups containing at least 12 carbon atoms, such as, for example, the poly-organosiloxanes described in patent application FR-A-2,641,185, and in particular polyorganosiloxanes containing a stearoyloxypropyl function;\nj) quaternary ammonium groups, such as in the products xe2x80x9cX2 81 08xe2x80x9d and xe2x80x9cX2 81 09xe2x80x9d and the product xe2x80x9cAbil K 3270xe2x80x9d from the company Goldschmidt;\nk) amphoteric or betaine groups, such as in the product sold by the company Goldschmidt under the name xe2x80x9cAbil B 9950xe2x80x9d;\nl) bisulphite groups, such as in the products sold by the company Goldschmidt under the names xe2x80x9cAbil S 201xe2x80x9d and xe2x80x9cAbil S 255xe2x80x9d.\nThe block copolymers having a polysiloxane-polyoxyalkylene linear block as repeating unit, which are used in the context of the present invention, preferably have the following general formula:\n([Y(R2SiO)aRxe2x80x22SiYO][CnH2nO)b])cxe2x80x83xe2x80x83(V)\nin which\nR and Rxe2x80x2, which may be identical or different, represent a monovalent hydrocarbon-based radical containing no aliphatic unsaturation,\nn is an integer ranging from 2 to 4,\na is an integer greater than or equal to 5, preferably between 5 and 200 and even more particularly between 5 and 100,\nb is an integer greater than or equal to 4, preferably between 4 and 200 and even more particularly between 5 and 100,\nc is an integer greater than or equal to 4, preferably between 4 and 1000 and even more particularly between 5 and 300,\nY represents a divalent organic group which is linked to the adjacent silicon atom via a carbon-silicon bond and to a polyoxyalkylene block via an oxygen atom,\nthe average molecular weight of each siloxane block is between about 400 and about 10,000, that of each polyoxyalkylene block being between about 300 and about 10,000,\nthe siloxane blocks represent from about 10% to about 95% of the weight of the block copolymer,\nthe average molecular weight of the block copolymer being at least 3000 and preferably between 5000 and 1,000,000 and even more particularly between 10,000 and 200,000.\nR and Rxe2x80x2 are preferably chosen from the group comprising alkyl radicals such as, for example, the methyl, ethyl, propyl, butyl, pentyl, hexyl, octyl, decyl and dodecyl radicals, aryl radicals such as, for example, phenyl and naphthyl, aralkyl radicals such as, for example, benzyl and phenethyl, and tolyl, xylyl and cyclohexyl radicals.\nY is preferably xe2x80x94Rxe2x80x3xe2x80x94, xe2x80x94Rxe2x80x3xe2x80x94COxe2x80x94, xe2x80x94Rxe2x80x3xe2x80x94NHCOxe2x80x94, xe2x80x94Rxe2x80x3xe2x80x94NHxe2x80x94COxe2x80x94NHxe2x80x94Rxe2x80x3xe2x80x94NHCO or xe2x80x94Rxe2x80x3xe2x80x94OCONHxe2x80x94Rxe2x80x2xe2x80x3xe2x80x94NHCOxe2x80x94, where Rxe2x80x3 is a divalent alkylene group such as, for example, ethylene, propylene or butylene, and Rxe2x80x2xe2x80x3is a divalent alkylene group or a divalent arylene group such as xe2x80x94C6H4, xe2x80x94C6H4C6H4xe2x80x94, C6H4xe2x80x94CH2xe2x80x94C6H4, C6H4xe2x80x94C(CH3)2xe2x80x94C6H4xe2x80x94.\nEven more preferably, Y represents a divalent alkylene radical, more particularly the xe2x80x94CH2xe2x80x94CH2xe2x80x94CH2xe2x80x94radical or the xe2x80x94C4H8xe2x80x94 radical.\nThe preparation of the block copolymers used in the context of the present invention is described in European application EP 0,492,657 A1, the teaching of which is included by way of reference in the present description.\nThe preferred polysiloxane-polyoxyalkylene linear block copolymers according to the invention are chosen from those of formula:\n[C4H8O(CnH2nO)b(CmH2mO)dxe2x80x94C4H8xe2x80x94SiMe2O(SiMe2)aSiMe2]cxe2x80x83xe2x80x83(VI)\nwhere Me represents methyl, n and m are integers ranging from 2 to 4, a is an integer greater than or equal to 4, or preferably between 5 and 200, b and d are integers greater than or equal to 0, preferably between 4 and 200, b+d is greater than or equal to 4, preferably between 4 and 200, and c is an integer greater than or equal to 4, preferably between 4 and 1000.\nAmong these copolymers, the ones more particularly used are those having a repeating unit of formula:\n[xe2x80x94(SiMe2O)xSiMe2xe2x80x94C4H8Oxe2x80x94(C2H4O)yxe2x80x94(C3H6O)zxe2x80x94C4H8xe2x80x94]xe2x80x83xe2x80x83(VII)\nwhere x is a number between 5 and 15 (limits included), y is a number between 15 and 30 (limits included), and z is a number between 20 and 40 (limits included).\nAmong these polymers, the ones more particularly used are those in which the siloxane/polyoxyalkylene weight ratio is about 20 and the polyoxyethylene/polyoxypropylene weight ratio is about 65/35.\nPolymers can also be chosen in which the repeating unit is of formula (VI) and whose siloxane/polyoxyalkylene weight ratio is about 75 and whose polyoxyethylene/polyoxypropylene weight ratio is about 50/50, polymers whose siloxane/polyoxyalkylene weight ratio is about 35 and whose polyoxyethylene/polyoxypropylene weight ratio is about 100/0, and polymers whose siloxane/polyoxyalkylene weight ratio is about 30 and whose polyoxyethylene/polyoxypropylene weight ratio is about 0/100.\nAccording to a particular embodiment of the invention, the block copolymer is chosen from the following copolymers:\n[[(CH3)2SiO]41(CH3)2SiCH2\nCH(CH3)CH2xe2x80x94O(C2H4O)18\nxe2x80x94(C3H6O)33CH2CH(CH3)CH2]\n16.1[[(CH3)2SiO]31\n(CH3)2SiCH2CH(CH3)CH2\nxe2x80x94O(C2H4O)20xe2x80x94(C3H6O)29\nCH2CH(CH3)CH2]13.3\n[[(CH3)2SiO]9(CH3)2\nSiCH2CH(CH3)CH2\nxe2x80x94O(C2H4O)20xe2x80x94(C3H6O)29\nCH2CH(CH3)CH2]26.3\n[[(CH3)2SiO]16(CH3)2\nSiCH2CH(CH3)CH2xe2x80x94O(C2\nH4O)18xe2x80x94(C3H6O)20CH2\nCH(CH3)CH2]21.5\n[[(CH3)2SiO]9(CH3)2\nSiCH2CH(CH3)CH2\nxe2x80x94O(C2H4O)5xe2x80x94CH2\nCH(CH3)CH2]4.8\nThe polymers containing a non-silicone organic skeleton grafted with monomers containing a polysiloxane, in accordance with the invention, are preferably chosen from those described in U.S. Pat. Nos. 4,963,935, 4,728,571 and 4,972,037 and patent applications EP-A-0,412,704, EP-A-0,412,707, EP-A-0,640,105 and WO 95/00578, the teachings of which are included in their entirety in the present description by way of non-limiting references. They are copolymers obtained by radical polymerization from monomers containing ethylenic unsaturation and from silicone macromonomers having a vinyl end group, or alternatively copolymers obtained by reaction of a polyolefin comprising functionalized groups and of a polysiloxane macromer having an end function which is reactive with the said functionalized groups.\nExamples of polymers containing a polysiloxane skeleton grafted with non-silicone organic monomers, which are suitable for carrying out the present invention, as well as the particular method for preparing them, are described in particular in patent applications EP-A-0,582,152, WO 93/23009 and WO 95/03776, the teachings of which are included in their entirety in the present description by way of non-limiting references.\nThe polyorganosiloxanes preferably used according to the invention are non-volatile polyorgano-polysiloxanes.\nThe polyorganosiloxanes are used in the compositions of the invention in proportions of between 0.01 and 50% by weight, and preferably between 0.1 and 30% by weight relative to the total weight of the composition.\nThe compositions according to the invention contain a cosmetically acceptable aqueous medium. They have a pH which can range from 3.5 to 11, preferably between 5.5 and 11 and even more preferably between 5.5 and 8.5.\nThe cosmetically acceptable medium for the compositions according to the invention consists more particularly of water and optionally of cosmetically acceptable organic solvents.\nThe organic solvents can represent from 0.5 to 90% of the total weight of the composition. They can be chosen from the group consisting of hydrophilic organic solvents, lipophilic organic solvents, amphiphilic solvents or mixtures thereof.\nAmong the hydrophilic organic solvents, mention may be made, for example, of linear or branched lower monoalcohols having from 1 to 8 carbon atoms, polyethylene glycols having from 6 to 80 ethylene oxide units, and polyols.\nAs amphiphilic organic solvents, mention may be made of polypropylene glycol (PPG) derivatives, such as esters of polypropylene glycol and of fatty acid, derivatives of PPG and of fatty alcohol, such as PPG-23 oleyl ether, and PPG-36 oleate.\nAs lipophilic organic solvents, mention may be made, for example, of fatty esters such as diisopropyl adipate, dioctyl adipate, alkyl benzoates and dioctyl malate.\nIn order for the cosmetic compositions of the invention to be more pleasant to use (softer when applied, more nourishing and more emollient), it is possible to add a fatty phase to the medium of these compositions.\nThe fatty phase can represent up to 50% of the total weight of the composition.\nThis fatty phase can contain an oil or a wax or mixtures thereof, and can also comprise fatty acids, fatty alcohols and fatty acid esters. The oils can be chosen from animal, plant, mineral or synthetic oils and in particular from liquid petroleum jelly, liquid paraffin, isoparaffins, poly-xcex1-olefins, fluoro oils and perfluoro oils. Similarly, the waxes can be chosen from animal, fossil, plant, mineral or synthetic waxes which are known per se.\nThe compositions of the invention can contain adjuvants that are common in the cosmetics field, such as other standard gelling agents and/or thickeners; emulsifiers; surfactants; moisturizers; emollients; sunscreens; hydrophilic or lipophilic active agents such as ceramides; anti-free-radical agents; sequestering agents; antioxidants; preserving agents; acidifying or basifying agents; fragrances; fillers; dyestuffs; reducing agents. The amounts of these various additives are those used conventionally in the fields considered.\nNeedless to say, a specialist will take care to select the optional compound(s) to be added to the composition according to the invention such that the advantageous properties intrinsically associated with the composition in accordance with the invention are not, or are not substantially, adversely affected by the addition envisaged.\nThe compositions according to the invention can be in any form which is suitable for topical application, in particular in the form of lotion-type solutions, in the form of aqueous or aqueous-alcoholic gels, in the form of vesicle dispersions or in the form of simple or complex emulsions (O/W, W/O, O/W/O or W/O/W emulsions) and can be of liquid, semi-liquid or solid consistency, such as milks, creams, gels, cream-gels, pastes and sticks, and can optionally be packaged as an aerosol and can be in the form of mousses or sprays. These compositions are prepared according to the usual methods.\nThe compositions according to the invention are preferably used as rinse-out or leave-in hair products, in particular to wash, dye, care for, condition, straighten or maintain the hairstyle or to permanently or temporarily reshape the hair.\nThe compositions can be styling products such as hairsetting lotions, blow-drying lotions, fixing compositions and styling compositions. The lotions can be packaged in various forms, in particular in vaporizers, pump-dispenser bottles or in aerosol containers in order to ensure application of the composition in vaporized form or in the form of a mousse. Such packaging forms are indicated, for example, when it is desired to obtain a spray or a mousse for fixing or treating the hair.\nThe compositions of the invention can also be shampoos, rinse-out compositions or leave-in compositions, to be applied before or after shampooing, dyeing, bleaching, permanent-waving or straightening the hair.\nThe compositions of the invention can also be used as hygiene or care products, such as protective, treatment or beauty creams for the face, for the hands or for the body, protective or care body milks, and skincare or skin cleansing lotions, gels or mousses.\nThe compositions of the invention can also be used as antisun compositions.\nThe compositions can also consist of solid preparations constituting cleansing soaps or bars.\nThe compositions of the invention can also be used as oral care products such as toothpastes and mouthwashes.\nThe compositions can be make-up products such as face creams, foundations, mascaras, eyeliners, lipsticks or nail varnishes.\nAnother subject of the invention is a cosmetic, non-therapeutic treatment process for the skin, the scalp, the hair, the eyelashes, the eyebrows, the nails or mucous membranes, characterized in that a composition as defined above is applied to the keratinous support, according to the usual technique for using this composition, for example application of creams, gels, sera, lotions or milks to the skin, the scalp or mucous membranes.\nThe examples which follow illustrate the invention without being limiting in nature."}
{"text": "The main components of wood are cellulose, hemicellulose and lignin. Cellulose and hemicellulose are hydrophilic structures that contain hydroxyl groups. Hydroxyl groups interact with water molecules to form hydrogen bonds. Consequently, wood is capable of absorbing as much as 100% of its weight in water producing swelling. Evaporation leads to shrinkage. Because this natural water absorption/evaporation cycle occurs non-uniformly, this cycle creates internal stresses within the wood. These stresses cause the wood to check, split and warp.\nResearch activities to improve the dimensional stability of wood have spanned many decades. Various approaches have been attempted to reduce the affinity of wood for water, such as heat treatment, chemical and enzymatic modification of hydroxyl groups on cellulose or hemicellulose, or provision of a barrier coating (either external or internal).\nCurrently, three commercial processes are available to impart dimensional stability to wood—thermal treatment, acetylation and furfurylation. Thermal treatment will improve the dimensional stability of wood, however, it also will cause significant loss of mechanical strength. In softwoods, acetylation generally confers an anti-swelling efficiency (ASE) of about 75% with an associated weight increase of about 26% to 28%. Acetylation requires impregnation of acetic anhydride into wood prior to initiation of acetylation. The acetylation process will generate by-product acetic acid in the treated wood. As a result, post-treatment removal of acetic acid is required; however, residual acetic acid will remain in wood. The generation of acetic acid used during wood treatment requires stainless steel, corrosion-resistant treating equipment. In addition, residual acid in the treated wood product requires the use of stainless steel and corrosion-resistant metal fastener hardware. Acetylated wood is also extremely prone to mold growth when exposed to a moist environment. As a result, acetylated wood requires a surface protective coating or colorants if used in an outdoor condition. Furfurylation generally provides treated wood with an ASE of about 60% and a weight gain of about 30%. Furfurylation processes and furfulated wood release undesirable volatile organic compounds (VOC) during the curing process. In addition, furfurylation also results in increased brittleness of the wood. The characteristics, expense and complexity of these processes for enhancing the dimensional stability of wood limit the commercial usefulness of these processes.\nResearch on improved treatments for enhancing the dimensional stability of wood has included cell-wall bulking treatments. The deposition of bulking agents in wood can be achieved by impregnating non-reactive bulking agents into the wood or by impregnating monomers into the wood followed by polymerization of the monomers within the wood. The bulking agents can be water soluble or insoluble, reactive or non-reactive with wood components. The bulking agents known to those skilled in the art include polyethylene glycol (PEG), phenol, resorcinol, melamine and urea-formaldehydes, phenol furfural, furfuryl-analine and furfuryl alcohol and various vinyl resins such as polystyrene, polymethyl methacrylate, polyacrylonitrile, polyvinyl chloride with the help of wood swelling agents. With the exception of PEG, most bulking agents penetrate into wood but remain in the cell lumen. As a result, these bulking treatments may temporarily retard water absorption by wood, but do not provide long-term stabilization, because the hydroxyl groups in the wood cell wall remain unmodified and are consequently still available to absorb water molecules.\nPolyether polyols are generally prepared by reacting epoxides, such as ethylene oxide, propylene oxide or tetrahydrofuran with initiator in the presence of catalyst(s). Common polyether diols are polyethylene glycol (PEG), polypropylene glycol (PPG), and poly(tetramethylene) glycol (also named polytetrahydrofuran or PTMEG). PEGs are generally water soluble, while PPG and PTMEG are water insoluble and can be pressure impregnated into wood.\nWater-soluble polyethylene glycol (PEG) polymers have been widely used for treating wood to maintain and improve the dimensional stability of wood and wood products by preventing shrinkage, drying and cracking, and reducing warping. In addition, pre-soaking green wood in a PEG solution can allow kiln drying of green wood at high temperatures without causing cracking, splitting or warping. PEG (even high-molecular weight PEG), however, is highly water soluble and renders it unsuitable for long-term outdoor use, because water and moisture will facilitate the leaching of PEG out of the wood.\nThe present inventors have discovered that water insoluble polyether polyols, such as polypropylene glycol (PPG) polymer or copolymer, or poly (tetramethylene) glycol (also named polytetrahydrofuran) polymer or copolymer, once impregnated into a wood product, provides excellent wood dimensional stabilization and is less prone to water leaching from the treated wood product than PEG."}
{"text": "As computer systems having high availability, HA clusters (High Availability clusters) and FT servers (Fault Tolerant servers) of a hot standby structure have been known.\nAn HA cluster is intended to make a system redundant by connecting a plurality of servers with one another. If a failure occurs in the currently used (active) server, processing is taken over by another server prepared as a standby system, so it seems that the cluster operates normally as a whole. Main methods thereof include an active-standby method and a replication method.\nIn an HA cluster of an active-standby method, the active system and the standby system share a storage. The active system writes application-dependent information necessary for synchronizing the standby system in the shared storage, and the standby system performs recovery processing with use of such information at the time of failover. As such, it is impossible to achieve availability transparently, when viewed from the applications and the OS. Further, it takes time for failover, and the services cannot be provided during that time.\nIn an HA cluster of a replication method, each of the active system and the standby system has a storage, independently. A request which reaches the application of the working system is transferred to the standby system, causing the standby system to have state transition which is the same as the active system. In general, making the states of a plurality of systems conform with each other is called synchronization. When a failure occurs in the working system and the working system stops, as the application state of the standby system is synchronized with the active system, it is possible to separate the active system and switch the processing to the standby system to thereby continue the service. However, as the replication mechanism must be added to each of the applications to be clustered, it is impossible to achieve the availability transparently when viewed from the applications and the OS.\nAs described above, in an HA cluster, it is necessary to add availability-conscious mechanism to the applications and the OS. On the other hand, in the case of an FT server, it is possible to continue services transparently without the need for particular processing by the applications and the OS. The approaches to realize an FT server include a hardware approach and a software approach.\nIn an FT server of a hardware approach, main hardware components such as a CPU, a memory, and a storage are made redundant, and if a failure occurs in any of the components, such a component is separated so as to continue operation. Under a definition that a module including a CPU, a memory, and a chip set is a CPU subsystem, and that a module including various IO devices is an IO subsystem, in a typical FT server in which components are duplexed, the methods for duplexing the CPU subsystem and for duplexing the IO subsystem are different. In the CPU subsystems, the operations of the hardware are caused to conform with each other completely in clock units. This is called lock-step. As the duplex CPU subsystems perform completely the same operations, when a failure occurs, the CPU subsystem in which the failure has occurred is separated logically, and the operation is instantly switched to the normal CPU subsystem to continue operation. In the IO subsystems, although they do not operate in lock-step, when a failure occurs, the operation is immediately switched to the other IO subsystem. An FT server of a hardware approach is able to realize an extremely high availability. However, as it is configured of special hardware, it takes a higher introduction cost compared with general servers of similar performance.\nOn the other hand, an FT server of a software approach uses a virtual technique which enables one or a plurality of OSs to operate on a physical computer. A computer virtually constructed on a physical computer is called a virtual computer or a virtual machine. In an FT server of a software approach, physical computers are made redundant, and a virtual computer of an active system and a virtual computer of a standby system are arranged on different physical computers, respectively. When a failure such as a hardware error occurs on the physical computer to which the virtual computer of the active system belongs, the processing performed by the virtual computer is continuously performed by the virtual computer of the standby system on the other physical computer. In order to continue the service transparently when viewed from the application and the OS, the FT server of the software approach performs processing to conform the states of the virtual computers of the active system and the standby system with each other, namely, synchronization.\nMethods for synchronizing the virtual computers of the active system and the standby system mainly include two methods, namely a virtual lock-step method and a checkpoint method. In a virtual lock-step method, an input to the virtual computer of the active system is also given to the virtual computer of the standby system so as to make the state of the virtual computer of the standby system transit in the same manner as that of the virtual computer of the active system. This method has an advantage that the quantity of data required for synchronization between the virtual computers is small. However, if the types of the CPUs of the active system and the standby system are different, there is a problem that the systems do not operate.\nOn the other hand, in a checkpoint method, images of the virtual computer of the active system (CPU, memory, storage, etc.) are periodically transmitted to the standby system so as to make the state of the virtual computer of the standby system conform with the state of the virtual computer of the active system. In the checkpoint method, implementation is relatively easy compared with the virtual lock-step method, and as it does not depend on a particular function of the CPU, there is an advantage that this method can be implemented in a variety of products. However, as images of the virtual computer have a large quantity of data, there is a problem that overhead for one time of synchronization is larger than that of the virtual lock-step method.\nIn order to solve this problem, a technique of transmitting only images of the virtual computer of the active system, which are updated after the previous checkpoint, has been proposed as first related art of the present invention (see Non-Patent Document 1 shown below, for example). In the first related art, when a checkpoint comes, the virtual computer of the active system is suspended so as to interrupt update to the main memory, and a local copy of all of the dirty pages, which are pages in the main memory having been updated after the previous checkpoint, is created in the buffer provided in the main memory. When a local copy has been created, the suspended virtual computer of the active system is restarted, and along with it, the copied dirty pages are transferred from the buffer to the standby system.    Non-Patent Document 1: Brendan Cully, and 5 others, “Remus: High Availability via Asynchronous Virtual Machine Replication”, [online], [searched on Sep. 5, 2012], Internet <URL: http://www.cs.ubc.ca/{tilde over ( )}andy/papers/remus-nsdi-final.pdf>\nHowever, creating local copies of all of the pages to be transferred in the memory requires a reasonable processing time. In the above-described first related art of the present invention, updating of the memory must be suspended until the entire data to be transferred has been copied completely. As such, the performance of the computer deteriorates."}
{"text": "Chinese culinary art is extensive and profound and has a variety of cooking techniques, which mainly comprise: stir frying, stewing, deep frying, boiling, steaming, pan frying, grilling, etc. The essence of various cooking techniques is to cook ever-changing dishes with various tastes and flavors by changing food match and controlling heating. Up to now, cooking process involves a lot of manual labor and relies very much on the cooker's experience and skills.\nThere exists some automated cooking machines, however, they normally cannot realize fully-automatic operation from feeding raw ingredients into the wok till pouring the cooked food out of the wok. More particularly, existing automated cooking machines cannot automatically pour the cooked food out of the wok or put accessory ingredients and raw food easily, and thus manual operation is still required during the cooking process. In addition, existing automatic operation program only has single function to cook dishes with same flavor and cannot be popular used."}
{"text": "This invention relates in general to a mechanical governor for an internal combustion engine, and more particularly to a cam shaft torque controlling device for an all-speed type of mechanical governor for a Diesel engine fuel injection pump in which the engine speed control range is set or determined by the position of a movable control lever.\nAs is well known, in existing governor assemblies, especially those for Diesel engines adapted to construction vehicles, speed control is implemented in the same manner for lower loads or lower torque ranges as it is for maximum load and torque conditions during a construction operation.\nFIG. 3 shows a plot of engine rotation speed N versus the position of the injection pump control rod R, which represents the amount of load exerted on or torque exerted by the engine. The rotation speed N.sub.B and the corresponding position of the control rod R.sub.B for a maximum engine output under full load setting is plotted, as is a predetermined low rotation speed N.sub.C, which is 50% lower than the rated speed, and the position of the control rod R.sub.C for a maximum engine output under full-load. In a conventional governor assembly, when the load is lower than R.sub.B speed control is effected as in an ordinal type of all-speed mechanical governor. If the load exceeds R.sub.B, however, the control rod must be shifted a distance Q (from R.sub.B to R.sub.C) in accordance with the speed deceleration curve from N.sub.B to N.sub.C in order to prevent the engine from suddenly stopping or stalling due to the increased (over capacity) load. This sharp speed slowdown alerts the operator to the overload, and enables him to back the machine off and/or shift to a lower gear ratio before the engine stalls. To implement this control function a torque spring is usually employed in the governor assembly. This presents a drawback, however, in that since the torque spring is operable or functioning at all times and for all load conditions, engine output shortages are unnecessarily caused thereby even during partial or reduced load operations."}
{"text": "This invention relates to a method and apparatus for producing waistband-equipped disposable diapers and, more particularly, to a method and apparatus wherein the usual processing steps performed or normal features added without having to go to extraordinary lengths to maintain the waistbands in their original tensioned condition. This has been a problem especially in diapers of the adult brief variety where the waistbands are much larger and stronger than in the infant variety.\nWaistband application in disposable diapers has been addressed in a number of prior art patents. Co-owned U.S. Pat. No. 4,523,969 proved successful in commercially producing adult diapers but with complicated machinery and processing. In the '969 patent and in U.S. Pat. No. 4,925,520 there were vacuum drums which were expensive to maintain.\nAnother prior art approach is seen in co-owned U.S. Pat. No. 5,000,806. This featured the diapers of the adult brief variety and which include a novel waistband construction and method of providing the same.\nThe prior art relating to waistbands can be seen in co-owned U.S. Pat. No. 5,000,806. This featured the application of an elastic strand manufactured by E. I. Du Pont de Nemours and Company marketed under the trademark Lycra. As Lycra can be stretched to levels exceeding 400%, it was difficult to obtain accurate placement of tapes or diecutting symmetry. This is overcome by the instant invention. Another prior art use of elastic waistbands is in U.S. Pat. No. 5,004,466 which employed a tape transfer step prior to adhesion to the poly-nonwoven substrates. The drawbacks of this method were high cost, unattractive tails visible through the material and the inability to control the placement.\nStill further, the prior art procedure was the use of a single strand elastic aligned on a drum-type former system. This left red lines due to concentrated elastic force to the wearer and the diaper had to be held in vacuum to remain aligned which resulted in a costly method as well as the elastic having to be rotated before application.\nThe shortcomings of the prior art are overcome by the instant invention which employs a center position nonwoven web to sandwich waistband material between the bottom moisture impervious sheet and the center nonwoven material. More particularly, the invention provides for applying spaced-apart longitudinally-extending webs of non-woven, performing the usual processing steps such as leg cutout and tape tab application, and then applying the waistbands. Other objects and advantages of the invention may be seen in the details of the ensuing specification."}
{"text": "As the typical process of fluorination for introducing fluorine atom into a substrate, heretofore, the direct fluorination using fluorine gas (F2) has been known (for example, refer to Patent Reference 1). When the substrate has a functional group such as oxygen, sulfur and halogen, processes in which the functional group is replaced with fluorine atom using an inorganic fluorinating agent such as hydrogen fluoride (occasionally referred to as HF, hereinafter) and sulfur tetrafluoride or an fluorinating agent other than inorganic fluorinating agents such as pyridine-9HF (the Olah reagent), a Yarovenko reagent of the fluoro-alkylamine type, a modified Ishikawa reagent of the fluoro-alkylamine type or diethylamionosulfur trifluoride (DAST), have been known (for example, refer to Non-Patent References 1 and 3).\nWhen the substrate is a halogen compound, the halogen-fluorine exchange reaction is the most convenient means for introducing fluorine. In the halogen-fluorine exchange, alkali metal salts of fluorine are used frequently. For example, sodium fluoride and potassium fluoride have little toxicity or possibility of causing corrosion unlike HF and can be handled easily. Potassium fluoride (occasionally referred to as KF, hereinafter) prepared in accordance with the spray drying process is frequently used recently (for example, refer to Non-Patent References 2 and 4).\nAs the agent other than those described above, molecular compounds of HF and a Lewis base such as pyridine or triethylamine, or ammonium fluoride salts can be used for the halogen-fluorine exchange reaction (for example, refer to Non-Patent References 3 and 4 (page 178)).\nThe processes for fluorination described above have a problem in that fluorine gas, hydrogen fluoride and sulfur tetrafluoride have toxicity and may cause corrosion and explosion, and special apparatuses and technologies are required for the handling.\nTo overcome the above problem and introduce fluorine safely and easily, various nucleophilic and electrophilic fluorinating agents have been developed (for example, refer to Patent References 2 and 3 and Non-Patent Reference 1). Patent Reference 3 relates to α,α-difluoro-amines proposed by the present inventors which are represented by the following general formula (1), can overcome the above problem of conventional processes, exhibit excellent heat stability and can be handled easily:\nwherein R0, R1 and R2 each represent hydrogen atom or an alkyl group, an aryl group, an alkylamino group or an arylamino group, which may have substituents, atoms and groups represented by R0, R1 and R2 may be same with or different from each other, and a ring may be formed by bonding of two or more groups represented by R0, R1 and R2.\nThe α,α-difluoroamine represented by general formula (1) can be produced in accordance with a conventional halogen-fluorine exchange reaction using as the precursor a halogen compound of an amide which is an α,α-dihaloamine represented by the following general formula (2):\nwherein R0, R1 and R2 each represent hydrogen atom or an alkyl group, an aryl group, an alkylamino group or an arylamino group, which may have substituents, atoms and groups represented by R0, R1 and R2 may be same with or different from each other, a ring may be formed by bonding of two or more groups represented by R0, R1 and R2, and X represents chlorine atom, bromine atom or iodine atom.\nThe α,α-difluoroamine represented by general formula (1) of the object compound can be obtained by the halogen-fluorine exchange reaction of the α,α-dihaloamine represented by general formula (2) using HF or an alkali metal salt of fluorine such as NaF and KF prepared in accordance with the spray drying process as the fluorine source.\nThe process for producing a fluorine compound by the fluorine exchange using HF or an alkali metal salt of fluorine such as NaF and KF has heretofore been known widely. The condition of the reaction can be decided with reference, for example, to Non-Patent Reference 3. However, occasionally, the reaction rate is insufficient, and it takes a long time to complete the reaction.\nFor example, when the reaction of N,N-diethyl-α-chlorometatoluoyl-amidium chloride is conducted using KF, which is prepared in accordance with the spray drying process, has a relatively great specific surface area and exhibits a great activity, in acetonitrile as the solvent under the refluxing condition (82° C.) for 24 hours, the yield is at most 70%. A long processing time not only lowers productivity of the object fluoroamines, but also results in an increased cost, and thus is a problem as an industrial process.\nAs described above, when an α,α-difluoroamine is produced in accordance with the halogen-fluorine exchange reaction using KF, the relatively great cost of KF prepared in accordance with the spray drying process and the long time required for the reaction are the major causes of the increase in the cost of production, and a further increase in the productivity and a further decrease in the cost are required from the standpoint of the industrial production.\nSome of the above problems can be overcome by using a molecular compound of HF and a Lewis base (occasionally referred to as “HF-Lewis base”, hereinafter) for the halogen-fluorine exchange reaction.\nHowever, few examples of the application of HF-Lewis base to the production of the α,α-difluoroamine represented general formula (1) can be found. In particular, no references clearly describe that triethylamine-3HF comprising HF and triethylamine in relative amounts by mole of 3:1 (causing no corrosion and enabling to use a glass vessel; occasionally referred to as “Et3N-3HF”) exhibits a greater nucleophilicity than that of other HF-Lewis bases such as pyridine-9HF (the Olah reagent) and is advantageously used for the halogen-fluorine exchange reaction since the reaction is rapidly completed.\nMore specifically, in conventional processes, no examples can be found on the production of N,N-diethyl-α,α-difluoro(3-methyl)benzylamine using HF-Lewis base such as Et3N-3HF in the chlorine-fluorine exchange reaction of N,N-diethyl-α-chlorometatoluoylamidium chloride. No examples can be found on the production of similar α,α-difluoroamines using HF-Lewis base, either.\nWhen the use, for example, of Et3N-3HF is applied to the production of α,α-difluoroamine, an unexpected difficulty arises as described in the following. When HF-Lewis base such as Et3N-3HF is used for the halogen-fluorine exchange reaction, the ratio of the amounts by mole of HF and the Lewis base is not always 1:1, and HF derived from HF-Lewis base used for the reaction is left remaining after the halogen exchange reaction is completed. This causes a serious problem in that separation of the product becomes difficult since molecular compounds are formed due to the interaction between HF and the nitrogen atom in the substrate or the product. Moreover, there is the possibility that the remaining HF causes corrosion. The process becomes complicated since steps for separation and purification must be added for obtaining the product.\nNo patents or references which mention the problems described above such as the problems in the industrial production of α,α-difluoroamine using Et3N-3HF or the means for solving the problems can be found.\n[Patent Reference 1] Japanese Patent Application Publication No. Showa 63 (1988)-25570\n[Patent Reference 2] Japanese Patent Application Laid-Open No. 2000-1477\n[Patent Reference 2] Japanese Patent Application Laid-Open No. 2003-64034\n[Non-Patent Reference 1] Yuki Gosei Kagaku Kyokaishi, 37, 1979, p. 606\n[Non-Patent Reference 2] Yuki Gosei Kagaku Kyokaishi, 47, 1989, p. 258\n[Non-Patent Reference 3] Journal of Organic Chemistry, 44, 1979, p. 3872\n[Non-Patent Reference 4] Chemistry of Organic Fluorine Compounds II, Monograph, American Chem. Soc., 1995, p. 187"}
{"text": "The present invention relates to selected metabolites of pirmenol, i.e., (+,-)-cis-.alpha.-[3-(2,6-dimethyl-1-piperidinyl)propyl]-.alpha.-phenyl-2- pyridinemethanol or N-[4-phenyl]-4-(2-pyridyl)-4-hydroxybutyl]-cis-2,6-dimethylpiperidine. The pharmacological activity of both pirmenol and the metabolites that are the present invention is antiarrhythmia. Thus, pharmaceutical compositions and use therefor comprising an antiarrhythmic effective amount of the metabolites are also within the scope of the invention.\nPirmenol is the starting material for synthetic preparation of the present invention metabolites. Therefore, U.S. Pat. No. 4,112,103, describing pirmenol, is incorporated by reference herein.\nPrecursors for pirmenol are disclosed in CA 97(11):92158m. These include .alpha.-[3-(2,6-dimethyl-1-piperidinyl)-1-propenyl]-.alpha.-phenyl[1(Z),2. alpha.,6.alpha.]-2-pyridinemethanol and .alpha.-[3-(2,6-dimethyl-1-piperidinyl)-1-propynyl]-.alpha.-phenyl[1(Z)-2. alpha.,6.alpha.]-2-pyridinemethanol. These disclosed precursors are different from the present invention because each precursor includes an unsaturated hydrocarbon chain. Also, the present compounds differ from both pirmenol and these precursors in the 2,6-dimethylpiperidine moiety."}
{"text": "1. Field of the Invention\nThe present invention relates to a zoom lens suitable for a television camera, a motion-picture camera, a video camera, a photography camera, and a digital camera, and more particularly, to a zoom lens having a high magnification, a small size and light weight, and little aberration deviation due to focus adjustment, and to an image pickup apparatus including the zoom lens.\n2. Description of the Related Art\nConventionally, as for a zoom lens in which focus adjustment is performed by a lens unit disposed closer to an object side with respect to a magnification-varying lens unit, there are proposed various types of systems.\nFor instance, Japanese Patent Application Laid-Open No. 2004-309761 discloses a zoom lens in which a focus lens unit having positive refractive power as a whole includes, in order from an object side, a first lens unit having negative refractive power, a second lens unit having negative refractive power, and a third lens unit having positive refractive power. In addition, Japanese Patent Application Laid-Open No. 2004-309761 proposes a method in which the second lens moves to the object side in focus adjustment to an object at a short distance.\nIn addition, Japanese Patent Application Laid-Open No. 2004-85846 discloses a zoom lens in which a focus lens unit having positive refractive power as a whole includes, in order from the object side, a first lens unit having positive refractive power, a second lens unit having negative refractive power, and a third lens unit having positive refractive power. Further, Japanese Patent Application Laid-Open No. 2004-85846 proposes a method in which the third lens unit moves to the object side in focus adjustment to an object at a short distance.\nThe zoom lens used for a television camera, a motion-picture camera, or the like is desired to achieve a higher magnification as well as a small size and light weight in order to secure mobility and to improve flexibility of photography. In addition, the zoom lens is required to have high performance with little aberration deviation due to focus adjustment.\nThe focus adjustment method of Japanese Patent Application Laid-Open No. 2004-309761 is suitable for a wide-angle zoom lens, but is difficult to achieve a high magnification. In the disclosed embodiment, a zoom ratio is approximately three. In order to achieve a high magnification, a diameter or an entire length of the first lens unit is increased.\nThe focus adjustment method of Japanese Patent Application Laid-Open No. 2004-85846 is suitable for a telephoto zoom lens, but is difficult to achieve a high magnification. In the disclosed embodiment, a zoom ratio is approximately 2.5. In order to achieve a high magnification, a diameter or an entire length of the first lens unit is increased."}
{"text": "The present invention relates in general to an adjustable supporting structure and, more particularly, to one which comprises adjustable flexible platform means capable of supporting heavy articles in different positions. Portions of the flexible platform means are carried by laterally spaced, longitudinal rolls the distance between which may be varied to vary the position in which a heavy article is supported.\nAn adjustable supporting structure of the foregoing general character is disclosed in Japanese Pat. No. 47-46421. More particularly, this patent discloses a slab turning device which includes two longitudinally spaced sets of flexible belts carried by laterally spaced, longitudinal rolls the distance between which may be varied to vary the position of the slab, the two sets of belts being driven by the rolls to turn the slab.\nThe present invention relates to an adjustable supporting structure of the foregoing general type, but one which can be used to support articles of different shapes and/or sizes in various positions so that any desired operations can be performed thereon. For example, the present invention may be utilized to support engines of different shapes and/or sizes in various positions for repair purposes, or the like, thereby providing a universal engine repair stand."}
{"text": "Prior analog to digital conversion techniques have used various methods. FIG. 1A represents a typical periodic sampling technique where an analog signal 10 is sampled every Tsample, so that the sampled values 12 may then be converted to a digital representation. FIG. 1B represents an alternative level-crossing technique where an input signal 14 is compared to a multiplicity of reference values 16, and the time periods, Dti, between reference level crossings of input signal 14 are recorded along with the specific level crossings."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a technique for rapid program compilation, and more particularly to a method for providing accurate floating-point computation.\n2. Description of the Related Art\nIn a floating-point computation in Java® (a registered trademark of Sun Microsystems Incorporated), precisions for a single precision computation and a double precision computation are defined by a language specification (24 bits mantissa and 8 bits exponent for the single precision computation, and 53 bits mantissa and 11 bits exponent for the double precision computation, respectively).\nWith CPUs of IA32 (Intel Architecture-32) architecture (the x86 family of Intel Corporation and compatible CPUs), it has been required to take either of the following approaches in order to execute both single precision computations and double precision computations:\n(Approach 1) switching the CPU's floating-point computation mode by setting the floating-point control word (hereinafter, referred to as FCW); or\n(Approach 2) executing the computation in the double precision mode, and if a computation result is to be obtained in single precision, storing the value on a floating-point register into a memory in single precision and reading it to degrade the computation precision.\nA method executing the single precision computation is considered by way of example. Generally, this method may be called in the double precision computation CPU mode (hereinafter, referred to as a double precision mode) or in the single precision computation CPU mode (hereinafter, referred to as a single precision mode). If the method is called in the single precision mode, it is more efficient to execute the code assuming the single precision mode. If the method is called in the double precision mode, it is necessary to prevent loss of computation precision according to Approach 1 or Approach 2 described above. One means of improving efficiency is to analyze the program to compare execution costs for the case using Approach 1 and for the case using Approach 2, and to execute the faster Approach, i.e. the one with less overhead (this analysis is hereinafter referred to as an inter-method analysis).\nIn a separate compilation environment for compiling a method having a high execution frequency, as in a dynamic compiler such as a JIT (Just In Time) compiler for Java®, it is not known whether the method being compiled will be called in the single precision mode or in the double precision mode. Therefore, for example, the precision mode is fixed as double precision at method boundaries (that is, all methods are called in the double precision mode), and if the single precision computation is executed, a code is generated according to the above described cost calculation.\nA technique for the above described cost calculation and code generation is described in the following prior art document, incorporated here by reference:\nM. Paleczny, C. Vick, and C. Click. The Java® HotSpot Server Compiler. The USENIX Association: Proceedings of the Java® Virtual Machine Research and Technology Symposium (JVM '01).\nThe above described prior art has a disadvantage in that the execution cost is high in either Approach 1 or Approach 2 and large overhead is incurred.\nSpecifically, in the case of Approach 1, a switch instruction for the FCW, with a form of, for example, (fldcw word ptr [mem]), is a high cost instruction. In addition, once the FCW is switched, it is necessary to switch the FCW again before and after a method call in that range and return the mode to the double precision mode as a standard, which may cause further overhead.\nIn the case of Approach 2, instructions for writing into and reading from a memory are executed whenever the single precision computation is executed; memory access may cause overhead and degrade overall execution performance.\nIn addition, in the dynamic compilation environment, an inter-method analysis may, for example, have excessive compilation time overhead, and so may not be effective for choosing a calculation Approach. Moreover, since new classes can be loaded dynamically, an analysis result may become invalid later, thus wasting computation resources.\nTherefore, it is an object of the present invention to provide an efficient mechanism for preventing loss of precision for both a single precision computation and a double precision computation in a floating-point computation, while avoiding execution overhead and invalidation of the analysis result."}
{"text": "With the continuous economic development and the growth of the financial industry, there are higher and higher requirements on the general performance, such as the quantity, the quality and the safety, of automatic teller machines applied in the financial industry. To adapt the rapid development of the financial industry, corresponding improvements are made to the automatic teller machines applied in the financial industry to satisfy the future development of the financial industry. However, a banknote picking up device of the automatic teller machine in the market generally has the following disadvantages.\nAt present, a banknote picking up mechanism and a banknote separating mechanism in the banknote picking up device of the automatic teller machine in the market are driven by a same banknote separating and picking up motor. In the case that a customer puts a stack of banknotes in the automatic teller machine but only wants to deposit 100 pieces of banknotes, when the hundredth banknote enters into a banknote separating wheel of the banknote separating mechanism, on one hand, if the banknote separating and picking up motor is stopped, the hundredth banknote will lose the power and stay in the banknote separating wheel, which will lead to banknotes loss and bring risk to the property of the customer; on the other hand, if the banknote separating and picking up motor still operates, extra banknotes will enter the banknote separating mechanism, thus a specific number of banknotes cannot be deposited, thereby the individual requirement of the customer cannot be satisfied.\nLikewise, as to other self-service equipment for processing sheet medium, it also cannot provide the service of depositing a specific number of sheet medium and suspending deposit for the customer.\nTherefore, there is an urgent need to provide a displaceable banknote picking up device being capable of depositing a specific number of banknotes to satisfy the individual requirement of the customer on one hand and suspending banknote deposit to ensure the safety of property on the other hand."}
{"text": "Ways to store objects that are frequently used and/or need to be easily accessible are known, but tend to require a good deal of installation hardware, work, and lack flexibility. For example, it is frequently desirous to store tools in a way in which they can be seen and easily accessed for use.\nTypically, a rack or pegboard is used, into which a set of hooks or is screwed. However, this storage solution requires a good deal of hardware and can be difficult to install. Further, once hooks or other tool fastening devices are screwed in or otherwise set, changing their configuration on the pegboard or rack to accommodate different tools or arrangement of tools can be difficult and take a good deal of effort. If the hooks or other holders are secured firmly, removing them to change their location may damage them or board, wall or rack they are secured to.\nVisible storage for toys is even more problematic. Currently, toys, such as, for example, toy guns that fire foam projectiles are typically stored in a box or other container. This can make them difficult to see or find. A child may forget what toys they have or at the least, have difficulty locating them. This can mean toys to not get used as they should, or even are purchased multiple times. Toys can also become damaged as children search containers for them, pull them out as they collide with other toys within the container, or even turn over the container, scattering toys, to locate a particular toy. If toys, tools or other objects that are accessed routinely, quickly or both are left out for easier access (or simply because there is no suitable place for them), the space around them can quickly become cluttered and annoying or difficult to navigate.\nWhat is needed is a system and method that addresses these issues, providing storage of items in a new and novel way that makes the items visible and easily accessible, yet is flexible enough to allow a user to easily change the configuration of storage."}
{"text": "Turbine shroud segments are exposed to hot gases and, thus, require cooling. Cooling air is typically bled off from the compressor section, thereby reducing the amount of energy that can be used for the primary purposed of proving trust. It is thus desirable to minimize the amount of air bleed of from other systems to perform cooling. Various methods of cooling the turbine shroud segments are currently in use and include impingement cooling through a baffle plate, convection cooling through long EDM holes and film cooling.\nAlthough each of these methods have proven adequate in most situations, advancements in gas turbine engines have resulted in increased temperatures and more extreme operating conditions for those parts exposed to the hot gas flow."}
{"text": "There exists a critical need to provide very large capacity rechargeable battery systems to the electrical utilities. This need comes about because the Nation's electrical generating capacity is frequently insufficient to provide peak power requirements. Storage batteries of the type described in this invention can be charged during off-power periods or during the night and then utilized to supplement the peak power demands. This application of very large capacity batteries is referred to in the electrical utility industry as load leveling or peak shaving. Batteries for this type of application must provide energy storage capacity at an economical price so as to be able to compete with alternate storage or peak power systems. As examples, compressed air and pumped water are alternate approaches to satisfying the peaking energy demand. It is generally accepted that a battery system with preferred features will utilize inexpensive constructional materials, will have very long-term cycling capability, will be environmentally harmless, will be inherently safe and because of real estate costs, will have a minimal size. Of great importance is the turnaround energy conversion efficiency which, as is generally accepted, should be 70 percent or greater.\nCertain efforts have already been made to develop systems which satisfy these requirements. The traditional lead-acid battery has been improved so as to be able to provide a minimum of 2500 cycles without deterioration. Even although these efforts have been partially successful, the cost figures of such batteries indicate that they will greatly exceed $30/kW-hr, which has been specified by the utility industry as the maximum allowable cost. Also, because of the nature of lead-acid batteries wherein the electrodes themselves are constantly being changed in composition over each charge-discharge cycle, it is difficult to accomplish very many cycles without some irreversible morphological change. To overcome this disadvantage, investigators have been developing a new type of battery referred to as a circulating electrolyte redox battery. In a redox battery the electrodes do not change in composition; they merely act as sites at which the electrochemical reduction and oxidation (redox) reactions of selected ions in solution can occur. Typical ion pairs may comprise titanium plus three ions and titanium plus four ions (Ti.sup.3+ /Ti.sup.4+) on one side of the battery and on the other, iron plus two ions and iron plus three ions (Fe.sup.2+ /Fe.sup.3+). In these redox cells, a corrosive acid supporting electrolyte is generally used, such as hydrochloric acid, and a separator which has ion selective properties is essential. The separator allows a common ion, such as the hydrogen ion, to transport the electrical charge from one side of the battery to the other while, at the same time, preventing the titanium ions and iron ions from mixing. If such mixing occurs, the cell becomes chemically short-circuited and the products of such mixing cannot be economically separated for reuse. Separators of the required degree of cationic discriminations suitable for this application have not yet been developed.\nThe open circuit voltage of redox cells is generally quite low, usually less than 1 volt, and during discharge this voltage is degraded by polarization. This polarization or degradation of potential is characteristic of many electrochemical systems, but in the redox system it increases with time as the ratios of the ion concentrations on each side of the battery change during discharge or charge. Because of the corrosive electrolyte, need for an ion selective membrane, and high chemical costs, it is unlikely that a redox system will satisfy the economic or durability requirements of electric utilities. Characteristics of redox load leveling battery systems are described in two recent publications; see, L. H. Thaller, Proceedings of the Symposium on Load Leveling, V 77-4, p. 353; N. P. Yao and J. R. Selman, Eds., The Electrochemical Society, Inc., Princeton, N.J., 1977; and J. Giner, L. Swette and K. Cahill, \"Screening of Redox Couples and Electrode Materials,\" Final Report, Contract NAS 3-19760, Sept. 1, 1976, CR 134705. In U.S. Pat. No. 3,553,017, sealed storage batteries based on selected pairs of redox couples are described. All of these couples require acidic electrolytes for proper functioning.\nAnother approach to development of load leveling batteries utilizes a zinc electrode as the negative electrode and a chlorine/chloride cell as the positive electrode. This system is illustratively described in U.S. Pat. Nos. 3,713,888; 3,935,024; 3,940,283; 3,954,502; and 4,001,036 and utilizes a circulating acid electrolyte. This zinc-chlorine system has several advantages: (1) mixing of the electrolyte adjacent to the electrodes does not degrade performance as happens with the aforementioned redox battery; (2) the voltage of the cell is quite high, i.e., 2.0 V; (3) the cell voltage does not change with state of charge or discharge because the chemical composition at the electrodes remains constant; and (4) relatively inexpensive reactants are utilized. Nevertheless, the system suffers from a number of disadvantages such as the corrosive nature of the acid electrolyte, self-corrosion of the zinc electrode in the acid electrolyte and generation of explosive hydrogen gas, polarization of the chlorine gas/chloride ion positive electrode, and the requirement for a refrigeration system to store the chlorine reactant as its hydrate in a low temperature aqueous solution."}
{"text": "The invention relates to communication between processors.\nMulti-processor computer systems have more than one processor. Each processor executes a separate stream (xe2x80x9cthreadxe2x80x9d) of instructions. It is sometimes necessary for two processors of a computer system to communicate data between themselves.\nIn one general aspect of the invention, a method of communicating between a first and a second processor includes the first processor sending a datum over a common control bus, and the second processor receiving the datum from the common control bus.\nAdvantages and other features of the invention will become apparent from the following description and from the claims."}
{"text": "1. Field\nThis application relates generally to wireless communication and more specifically, but not exclusively, to controlling access point transmissions.\n2. Introduction\nA wireless communication network may be deployed over a geographical area to provide various types of services (e.g., voice, data, multimedia services, etc.) to users within that geographical area. In a typical implementation, macro access points (e.g., each of which provides service via one or more cells) are distributed throughout a macro network to provide wireless connectivity for access terminals (e.g., cell phones) that are operating within the geographical area served by the macro network.\nAs the demand for high-rate and multimedia data services rapidly grows, there lies a challenge to implement efficient and robust communication systems with enhanced performance. To supplement conventional network access points (e.g., to provide extended network coverage), small-coverage access points (e.g., low power access points) may be deployed to provide more robust indoor wireless coverage or other coverage to access terminals inside homes, enterprise locations (e.g., offices), or other locations. Such small-coverage access points may be referred to as, for example, femto cells, femto access points, home NodeBs, home eNodeBs, or access point base stations. Typically, such small-coverage access points are connected to the Internet and the mobile operator's network via a DSL router or a cable modem. For convenience, small-coverage access points may be referred to as femto cells or femto access points in the discussion that follows.\nWhen a femto cell is deployed on a carrier frequency that is different from the carrier frequencies used by neighboring macro cells, the femto cell may radiate beacons on the macro cell carrier frequencies. In this way, the femto cell may attract an access terminal that is in the vicinity of the femto cell to the femto cell coverage (i.e., cause the access terminal to move off of the macro cell coverage). Thus, through the use of this beacon scheme, a user coming home (e.g., approaching a home femto cell) from outside the coverage of the femto cell will be able to readily discover the femto cell and obtain service from the femto cell. Though such beacons are useful in terms of femto cell discovery, they may create interference on the macro network since the beacons are transmitted on the same carrier frequency that is used by neighboring macro cells. This interference may affect the voice call quality of active macro cell users (i.e., users actively receiving service from one or more macro cells on a macro cell frequency) and may also lead to call drops if the macro cell user happens to be very close to the femto cell. Therefore, there is a need to protect active macro cell users from interference from femto cells."}
{"text": "Recently, Scanning Electron Microscopy (hereinafter referred to as “SEM”), Transmission Electron Microscopy (hereinafter referred to as “TEM”), or Scanning Transmission Electron Microscopy (hereinafter referred to as “STEM”) are often used to analyze a sample having a fine structure. Before introducing a sample into these observation devices, it is necessary to prepare a cross section or thin film of the sample, and thus a Focused Ion Beam (hereinafter referred to as “FIB”) device is used. A method for FIB processing an observation portion into a thin film shape is called FIB microsampling, and this is the most suitable sample preparation method when using electron microscopy and the like to analyze the condition/structure on the order of several nm, which is the subject of recent nanotechnology research (PTL 1).\nIn processing or observation of a sample using a charged particle beam device, the temperature of the sample rises due to the energy of the charged particle beam, and there are cases in which it can be difficult to analyze the original structure of the sample. Thus, a method of processing or observing a sample using a charged particle beam device while cooling the sample has been proposed (PTL 2).\nFurther, in the procedure for preparing a thin film sample in order to observe the cooled sample with a TEM device or STEM device, a method of using a cooled manipulator to remount a sample with a thin film shape prepared on a base material onto a cooling sample holder has been proposed as disclosed in PTL 4.\nIn addition, among methods for preparing a cross section of a material that is sensitive to heat such as a liquid sample, a biological sample, or a polymer material without using a FIB device, a method of cutting a frozen sample with a cooled knife such as freeze etching is known (NPL 1)."}
{"text": "Glass substrates may be coated, for example, to provide a decorative effect or to enhance substrate properties. For instance, glass substrates may be coated to provide anti-shattering properties, abrasion resistance, increased elasticity, and solvent resistance. Glass containers, for example, may benefit from clear coatings that provide mechanical protection to the external surfaces to help minimize mechanical damage, such as scuffing or marring, to the containers during transportation, storage, filling operations, and distribution."}
{"text": "This invention relates to apparatus which uses flowable nanoparticles for printing images.\nInk jet printheads are used to selectively eject ink droplets onto a receiver to form an image. Within the printhead, the ink may be contained in a plurality of channel members and energy pulses (either piezo or heat) are used to actuate the printhead channel members or the ink orifices causing the droplets, which form the reservoirs, of ink to be ejected on demand or continuously, through an orifice plate which is placed over the channel members.\nIn one representative configuration, a piezoelectric ink jet printing system includes a body of piezoelectric material defining an array of parallel open topped channel members separated by walls. In the typical case of such an array, the channel members are micro-sized and are arranged such that the spacing between the adjacent channel members is relatively small. The channel walls have metal electrodes on opposite sides thereof to form shear mode actuators for causing droplets to expel from the channel members. An orifice defining structure includes at least one orifice plate defining the orifice through which the ink droplets are ejected, and is bonded to the open end of the channel members. In operation of piezoelectric printheads, ink is directed to and resides in the channel members until selectively ejected therefrom. To eject an ink droplet through one of the selected orifices, the electrodes on the two side wall portions of the channel in operative relationship with the selected orifice are electrically energized causing the side walls of the channel to deflect into the channel and return to their normal undeflected positions when the applied voltage is withdrawn. The driven inward deflection of the opposite channel wall portions reduces the effective volume of the channel thereby increasing the pressure of the ink confined within the channel to force few ink droplets, 1 to 100 pico-liters in volume, outwardly through the orifice. Piezoelectric ink jet printheads are described in detail in U.S. Pat. Nos. 5,598,196; 5,311,218; 5,365,645, 5,688,391, 5,600,357, and 5,248,998. Alternative ink jet print head configuration utilizes thermal energy to eject ink droplets from the orifices onto the receiver. Thermally activated ink jet print heads are described in details in U.S. Pat. Nos. 4,849,774; 4,500,895; and 4,794,409. This process of forming channel members, particularly in piezoelectric materials, is not only time consuming and expensive, but also is amenable to many defects generated during cutting the channel members or forming the channel members thereby reducing the throughput and increasing the unit manufacturing cost. Furthermore, mechanical damages caused during sawing or laser cutting also are detrimental to the piezoelectric characteristics of the material.\nAnother significant problem encountered in ink jet printing is the drying of the ink either inside the channels or at the orifice plates or at the orifices. To overcome this problem the ink formulators routinely compromise in ink formulations and the system designers incorporate wiper blades at the orifices or depend on ultrasonic devices to clean and remove the dried ink. Again, this introduces additional complicated system architecture and consequently increases the unit manufacturing cost.\nInks provide their own set of problems in ink jet printing, they are at all times flowable, thus subject to leakage, and may have storage problems, including limited shelf life. They also have a problem of drying and difficulty maintaining uniform viscosity since they can be temperature sensitive. Inks have been formed which have numerous ingredients such as dyes, pigments and colorants and nanoparticles (see U.S. Pat. No. 5,679,138). But, in all cases, the ink remains flowable in all states and subject to problems.\nIt is an object of the present invention to proved an improved apparatus which makes use of material having nanoparticles which is not a flowable ink for forming images on receivers.\nThis object is achieved in an apparatus for printing an image using image forming materials having nanoparticles less than an average diameter of 500 nm, the apparatus defining a plurality of printing nozzles each of which prints a pixel on a receiver, the nanoparticles being effective in two states, in a first state it aggregates and will not flow through an ink jet printer nozzle and in a second state when subject to a force is flowable so as to be deposited on the receiver, such nanoparticles being bound to each other by Van der Waals forces in the first state and after deposition on the receiver comprising:\n(a) a hopper for receiving the nanoparticles and defining a cavity into which the nanoparticles aggregate in the first state; and\n(b) means for applying a force to the nanoparticles disposed in the hopper to cause them to be effective in the second state and deliver them to the nozzles where such nanoparticles flow and are deposited on the receiver.\nIt has been discovered that nanoparticles which are effective in two states can be used to form images on receivers. Nanoparticles having an average diameter of less than 500 nm aggregate because of Van der Waals forces and are not flowable past an orifice in a first state but are flowable in a second state when a force as been applied to overcome the Van der Waals forces. A feature of the invention is that a stream of gaseous material, such as air, can be directed so as to cause the flow of such nanoparticles towards the receiver.\nThe nanoparticle imaging materials used with this system have no vapor pressure and hence will not dry out and clog, eliminating the need for such complexity in the print head. The present invention provides a simple solution to ink clogging problems and facilitates the use of smaller orifices. Smaller orifices can provide greater resolution. Additionally, the speed of the ink jet printer system can be increased because there is no fluid component to the nanoparticle imaging materials which eliminates the need to wait for solvent drying on the receiver, a major limitation in ink jet printers."}
{"text": "The replacement plate of the invention replaces the damaged lock opening in the trunk of a vehicle. The replacement plate of the invention restores a car trunk to its original appearance. The plate is affixed to the body of the vehicle over the opening left by the previous lock, which was removed or broken out, so that a new replacement lock may be installed. This restores the trunk to a lockable condition. The keyway slots of the protruding center of the escutcheon plate permit the mounting of a replacement lock with facility and rapidity."}
{"text": "In recent years, wireless communication systems have emerged as an important alternative to wireline communications in the United States. These systems offer users the convenience of mobility and personal services, as well as significant cost savings over traditional wireline systems in the installation and relocation of telecommunications equipment. The use of wireless systems is also proliferating in remote areas of the world due to the ease of installation of such systems as compared to wired communications.\nOne type of wireless communications system, cellular telephone systems, uses radio communications between a plurality of subscriber units with in the wireless system and between subscriber units and the Public Switched Telephone Network (\"PSTN\") for calls to or from outside the wireless system.\nIn cellular systems, as illustrated in FIG. 1, a geographical area is divided into cells covered by the radio system at the base station. The cells are connected via a respective base station (B1-B7) to a Mobile Services Switching Center (\"MSC\"). The MSC, in turn, is connected to the PSTN. The telephone users (mobile units) communicate with the MSC through the base stations. The MSC switches calls between wireline and mobile subscribers, controls signaling to the mobile stations, compiles billing statistics, and controls the overall operation of the system.\nCellular systems overcome many of the problems associated with traditional wireline communications. These systems, however, still require the installation of expensive infrastructure equipment (i.e., base stations B1-B7, and MSC's) to provide system functionality and the necessary switching capabilities.\nOther wireless communications systems overcome the need for expensive infrastructure equipment by using point-to-point radio communications. However, these systems do not provide the extensive switching capabilities provided in cellular and PCS systems. For example, multiple-user radio transceiver communication systems exist in which squelching circuits are used to limit communications between selected radio transceivers within a group of transceivers. In one such system described in U.S. Pat. No. 4,646,358, a first transceiver sends an origination address and a destination address. The transceivers within the broadcast area receive these addresses. When the received destination address corresponds to the address of the second transceiver, the second transceiver turns off its squelch circuit in order to receive the communications from the first transceiver. Any transceiver that receives the radio broadcast and does not have a corresponding destination address leaves its squelch circuit on so that the received signal cannot be heard by the user.\nThe second transceiver then communicates with the first transceiver by sending its address and the origination address of the first transceiver with a broadcast communication. When the first transceiver receives the broadcast communication with its address, the first transceiver turns off its squelching circuit in order to receive the communications. Thus, two transceivers may establish communications in a broadcast radio communications system. In operation, the system requires that all signaling between transceivers include both the address code of the originating station and the address code of the destination station. However, this system suffers from several disadvantages. For example, because origination and destination addresses must be sent with all communications, the likelihood of coding and transmission errors increases. Therefore, if an origination address is incorrectly received, stored, or retransmitted, the intended transceiver may not receive the communication.\nIn another system using squelching circuits, as described in U.S. Pat. No. 5,058,205, the origination transceiver only transmits the address with the initial message transmission. Once communication is established between two transceivers, a responding transceiver does not send its own address nor does the originating transceiver send its own address again during two-way communications between the two transceivers. Instead, the only address information required to be sent between the transceivers is the destination address of the communications.\nThese point-to-point radio communication systems, however, do not provide the capabilities associated with traditional wireless systems. In particular, these systems do not permit these transceivers to send and receive calls via the PSTN, maintain billing information, or provide switching capabilities associated with cellular telephone systems.\nThus, a need has arisen for an improved wireless communication system that eliminates the problems of existing systems, while maintaining the convenience of mobility and personal services provided by traditional wireless systems."}
{"text": "The present invention relates generally to refrigerators, and, more specifically, to icemakers therein.\nA typical refrigerator includes a cool compartment for keeping cool frozen food and beverage products, and a freezer compartment for storing frozen products.\nA common feature located in modern refrigerators is an automatic icemaker built into the freezer compartment. The icemaker includes an ice tray having several compartments in which individual ice cubes are formed. The frozen ice cubes are mechanically dispensed from the ice tray into an ice bin from which the ice may be conveniently dispensed through a chute and outlet typically formed in the freezer door.\nThe empty ice tray is then refilled with water for the next batch of ice cubes. This is automatically done by the timed opening of a solenoid-actuated valve.\nHowever, accurately controlling the volume of water delivered to the ice tray by timing water flow is dependent upon water pressure. Timed filling of the ice tray must therefore be selected to accommodate maximum expected water pressure without overfilling the ice tray, which will result in underfilling the ice tray for water pressure levels therebelow.\nFurthermore, competitive pricing of icemakers requires simplicity of design, few components, and compact elements for minimizing space requirements within the cabinet of the refrigerator. And, electrically operated components, such as solenoids, should be configured for reducing heat generation therefrom which adversely affects performance of the refrigerator.\nAccordingly, it is desired to provide an improved dispenser for accurately filling icemaker trays without regard to variation in water pressure."}
{"text": "1. Field of the Invention\nThis invention relates to a device to secure articles to a bag, hook, item of clothing or other item. More particularly, this invention relates to a holder having a hook or loop type fastening surface for retaining articles, and a snap-hook type fastener for securing the holder to a golf bag or other item.\n2. Description of the Related Art\nNumerous means for holding gloves and other articles to a golf bag or other item are known. For example, U.S. Pat. No. 5,177,881 discloses a device that holds a golf glove using Velcro® type (hook and loop type) fasteners. Two straps are attached to either side of a decorative plate-like element, with one strap holding the glove via hook and loop type fasteners and the other strap attached to the golf bag via a hook type fastener. The device can be hung from a golf bag.\nU.S. Pat. No. 5,864,925 discloses a device that holds a golf glove to a belt or other item using a plastic card having a patch of hook and loop type fasteners. The glove is attached to the patch of hook and loop fasteners, and the card is suspended from a belt via a metal clip.\nU.S. Pat. No. 6,546,599 discloses a device for carrying a golf glove from a back pocket. The device is a foldable rectangular panel having a surface partially covered with hook and loop type fastening material. The panel can be folded and one half inserted into the back pocket of the user while the part covered with hook and loop type fastening material hangs outside the back pocket. Magnets may be used to secure the half of the panel inside the pocket to the half outside the pocket.\nU.S. Pat. No. 6,694,575 discloses a flat, rectangular sports holder that attaches to a belt. The holder includes a patch of facing hook and loop type fastening material for securing the glove to the holder and slits threw which the belt is threaded.\nWhile each of these inventions may be suitable for its intended purpose, none of these inventions provides a holder having two rigid halves hingedly attached by a connecting strap such that the holder can be folded and unfolded to hold not only items having a hook and loop type fastening surface, but also items having straps or other means for draping the item over the connecting strap. It is an object of the present invention to provide a holder having these features. Another object of the present invention is to provide a device that can hold a glove or other item and is easy to use and attractive in appearance. Further and additional objects will appear from the description, accompanying drawings, and appended claims."}
{"text": "Options for accessing and listening to digital audio in an out-loud setting were limited until in 2003, when SONOS, Inc. filed for one of its first patent applications, entitled “Method for Synchronizing Audio Playback between Multiple Networked Devices,” and began offering a media playback system for sale in 2005. The Sonos Wireless HiFi System enables people to experience music from many sources via one or more networked playback devices. Through a software control application installed on a smartphone, tablet, or computer, one can play what he or she wants in any room that has a networked playback device. Additionally, using the controller, for example, different songs can be streamed to each room with a playback device, rooms can be grouped together for synchronous playback, or the same song can be heard in all rooms synchronously.\nGiven the ever growing interest in digital media, there continues to be a need to develop consumer-accessible technologies to further enhance the listening experience.\nThe drawings are for the purpose of illustrating example embodiments, but it is understood that the inventions are not limited to the arrangements and instrumentality shown in the drawings."}
{"text": "Hydrocracking refers to a process in which hydrocarbons crack in the presence of hydrogen and catalyst to lower molecular weight hydrocarbons. Depending on the desired output, the hydrocracking unit may contain one or more beds of the same or different catalyst. Hydrocracking is a process used to crack hydrocarbon feeds such as vacuum gas oil (VGO) to diesel including kerosene and gasoline motor fuels.\nMild hydrocracking (MHC) is generally used upstream of a fluid catalytic cracking (FCC) or other process unit to improve the quality of an unconverted oil that can be fed to the downstream conversion unit, while converting part of the feed to lighter products such as diesel. As world demand for diesel motor fuel is growing relative to gasoline motor fuel, mild hydrocracking is being considered for biasing the product slate in favor of diesel at the expense of gasoline. Mild hydrocracking may be operated with less severity than partial or full conversion hydrocracking to balance production of diesel with the FCC unit, which primarily is used to make naphtha. Partial or full conversion hydrocracking is used to produce diesel with less yield of the unconverted oil which can be fed to a downstream conversion unit.\nDue to environmental concerns and newly enacted rules and regulations, saleable diesel must meet lower and lower limits on contaminates, such as sulfur and nitrogen. New regulations require essentially complete removal of sulfur from diesel. For example, the ultra low sulfur diesel (ULSD) requirement is typically less than 10 wppm sulfur.\nHydrotreating refers to a process in which olefins and aromatics are saturated and heteroatoms, such as sulfur, nitrogen and metals are removed from the hydrocarbon feedstock over catalyst in the presence of hydrogen. Hydrotreating is an essential step in the production of ULSD.\nHydroprocessed streams are typically stripped with an inert gas such as steam to remove volatile ammonia and hydrogen sulfide to reduce sulfur and nitrogen concentration in the product fuel stream. Stripped hydroprocessed streams are then fractionated to provide product fuel streams.\nMHC reactors are typically operated at low to moderate conversion and lower pressures than higher conversion hydrocrackers, so the distillate produced from MHC units can be high in sulfur, such as 20 to 150 wppm sulfur, because the environment in the MHC reactor has a high concentration of hydrogen sulfide. In addition, the high concentration of ammonia in the MHC reactor reduces hydrocracking activity requiring higher operating temperatures further limiting organosulfur conversion. As a result, diesel from the MHC reactor must be also treated in a hydrotreating unit to achieve ULSD.\nThere is a continuing need for improved methods of producing more diesel from hydrocarbon feedstocks than gasoline. Such methods must ensure that the diesel product meets increasingly stringent product requirements."}
{"text": "It is necessary that an anchoring device settled in the ground be capable of bearing a pulling load as large as possible when it is adapted for supporting a standing object.\nTherefore, it is an object of the present invention to provide a novel anchoring device capable of exhibiting an increased resistance to pulling loads exerted on the device on the surface of the ground, and a method for rigidly settling the device in the bottom of a deep shaft in the ground in a manner such that the settled device cannot be pulled out of the shaft.\nIt is another object of the present invention to provide a novel anchoring device provided with means for ascertaining that the device is appropriately settled in the ground.\nIt is still another object of the present invention to provide an anchoring device provided with means for bearing compression exerted on the device, and a method for settling the device in the ground so that the anchoring device can bear the compression applied from above the ground.\nIn accordance with one of the features of the present invention, an anchoring device comprises\nA VERTICAL TUBULAR ROD; A PLURALITY OF SUPPORT ARMS, ONE END OF EACH SUPPORT ARM BEING CONNECTED WITH A LOWER PART OF THE TUBULAR ROD SO THAT SAID SUPPORT ARM IS TURNABLE ABOUT THE CONNECTED END IN A VERTICAL PLANE INCLUDING AN AXIS OF SAID TUBULAR ROD; A PLURALITY OF FLUKE BLADES ONE OF WHICH IS MOUNTED ON EACH SUPPORT ARM, TO CONSTITUTE CUTTING ARMS; MEANS FOR URGING SAID CUTTING ARMS TO OUTWARDLY TURN ABOUT SAID CONNECTED ENDS OF SAID SUPPORT ARMS FROM A FOLDED STATE IN WHICH SAID CUTTING ARMS ARE BROUGHT TOGETHER ABOUT SAID TUBULAR RODS; MEANS FOR STATIONARILY SUPPORTING ONE END OF SAID URGING MEANS ON A PART OF SAID TUBULAR ROD; HOLDING MEANS FOR RELEASABLY RESTRAINING SAID CUTTING ARMS IN THE FOLDED STATE, SAID HOLDING MEANS HAVING A LINE MEANS TO RELEASE THE RESTRAINT OF SAID CUTTING ARMS ABOVE THE GROUND SURFACE, AND; MEANS FOR LIMITING THE TURN OF SAID CUTTING ARMS TO A SELECTED AMOUNT OF ANGLES FROM SAID FOLDED STATE.\nFurther, in accordance with another feature of the present invention, a method for settling the anchoring device of the present invention into a shaft bored in the ground comprises at least the processes of:\nI. FOLDING THE CUTTING ARMS BY MEANS OF THE HOLDING MEANS HAVING THE RELEASE LINE;\nII. PULLING SAID RELEASE LINE OF SAID HOLDING MEANS WHEN SAID ANCHORING DEVICE REACHES A BOTTOM OF SAID SHAFT;\nIII. PROVIDING ROTATIONAL MOTION TO SAID CUTTING ARMS LYING AGAINST THE INNER WALL OF SAID SHAFT FROM ABOVE THE GROUND BY MEANS OF AN ORDINARY BORING MACHINE FOR A PREDETERMINED TIME, AND;\nIV. LIFTING SAID ANCHORING DEVICE WHILE SAID CUTTING ARMS ROTATE AT A PRESELECTED HEIGHT SO THAT A CAVITY ENGAGEABLE WITH SAID CUTTING ARMS OF SAID ANCHORING DEVICE IS FORMED IN SAID BOTTOM OF SAID SHAFT.\nThe other features, and advantages of the present invention will be apparent from the ensuing description with reference to the accompanying drawings."}
{"text": "This invention relates to battery plates and method for making same and, more particularly, to rechargeable battery plates in which the active material is impregnated in a sintered porous plaque.\nIn U.S. Pat. No. 3,790,408 assigned to the assignee of this invention, there is described a battery coil in which one of the longitudinal edges of the negative plate is extended transversely beyond one longitudinal edge of the positive plate at one end of the coil and the other longitudinal edge of the positive plate is extended beyond the other longitudinal edge of the negative plate at the opposite end of the coil. This construction permits the easy attachment of positive and negative plate current collectors to the extended edges of the respective plates by welding of the current collectors to the extended edges at a plurality of contact points. The plates are comprised of a thin metal substrate with a sintered porous plaque coating both major surfaces of the substrate. The active material of the battery is impregnated in the pores of the porous plaque. A small region adjacent one longitudinal edge of each plate which forms the extended edge is free of the porous plaque and active material for, among other reasons, assuring the formation of a good weld between the plate edge and the current collector. At the opposite longitudinal edge of each plate, the porous plaque extends to the substrate edge and terminates in a flat edge coplanar with the substrate edge (hereinafter referred to as \"flat edge termination\"). The principal reasons for flat edge termination is that it is important to maximize the volume of porous plaque per unit cell volume in order to maximize the energy density of the cell. This is particularly important in small rechargeable cells such as nickel-cadmium.\nIn the method for making such plates, an elongated metal substrate is used which has a width dimension such that a plurality of strips of plates can be formed by longitudinally severing the substrate into a plurality of strips. The substrate is coated on both major surfaces with a slurry comprised of a metal powder. The slurry coating is then divided into continuous strips by forming longitudinally slurry-free channels by using doctor blades to wipe the channels free of slurry. Also, a region adjacent each edge of the substrate is wiped free of slurry. The edge regions of the substrate eventually form the extended edge of each plate which is to be formed adjacent thereto.\nEach continuous strip of slurry is dimensioned to be approximately equal either to the sume of the widths of two battery plates or to the width of a single battery plate (hereinafter referred to as \"double width\" and \"single width\" strips). For example, when using a substrate of 20.32 cm width, three double width strips of slurry are defined by two channels by running longitudinally of the substrate. The double width strips eventually form plates of 3.2 cm width used in a sub-C cell. A variation to accommodate plates of other widths is to form, on a substrate of 18.0 cm width, one double width strip and one songle width strip. The double width and single width strips are used to form plates of 4.7 cm used in D cell. Other variations are also possible. After the slurry strips are formed, the substrate is passed through a sintering furnace in which the slurry is transformed into a porous plaque. The porous plaque coated substrate is then severed into strips of plate of the finished width. The strips of plate cut from the double width strips have a slurry-free edge region along one longitudinal edge which is formed from the slurry-free channel or edge region of the substrate and the opposing longitudinal edge of the plate strip has a flat edge termination. Similarly, the single width strips are cut so that one longitudinal edge has a slurry-free edge region formed by the slurry-free edge of the substrate and the opposing longitudinal edge has a flat edge termination. The strips are then impregnated with active material and transversely severed into a plurality of plates.\nWhile this construction has been found to be generally advantageous, it has been noted that there is a tendency for the cells to become shorted during cell assembly and testing. Cell shorting is also one of the major problems which will cut short the expected life of such cells.\nAccordingly, it is an object of this invention to provide an improved porous plaque plate construction which is less subject to shorting during assembly into battery coils and during their subsequent use as cells.\nAnother object of the invention is to provide a porous plaque battery plate which is less subject to flaking and chipping of the porous plaque during assembly."}
{"text": "Clothes washing machines can either be vertical or horizontal axis for moving clothes during the washing operation. In the horizontal washing machine, the drum or basket rotates either in one direction or intermittently in both directions causing the clothes therein to tumble during the washing operation in the soap and water cleaning solution. In the vertical washing machine, the agitator reciprocates or oscillates to continually change its rotation direction for moving the clothes to effect cleaning thereof in the soap and water solution.\nIn order to improve clothes moving in the vertical washing machine, it is known to also include an auger having a spiraling vane or screw disposed coaxially atop the agitator. Disposed between the auger and the agitator is a conventional ratchet mechanism which converts the oscillating, bidirectional rotation of the agitator into unidirectional rotation of the auger so that the auger screw is rotated for pulling the clothes downwardly in operation in the cleaning solution for improving the effectiveness of the washing operation.\nA typical ratchet mechanism includes a disk attached to the top of the agitator from which extends radially outwardly therefrom a plurality of ratchet pawls. A cooperating ratchet wheel in the form of a ring having a plurality of radially inwardly facing ratchet teeth attached to the bottom of the auger is disposed in a common axial (horizontal) plane with the pawls. Oscillation of the agitator in turn imparts force on the pawls which intermittently are driven radially outwardly for engagement with the cooperating ratchet teeth. As the agitator rotates in one direction, the pawls are driven radially outwardly and engage the ratchet teeth for also rotating the auger in the same direction. However, when the agitator rotates in the opposite direction, the pawls disengage the ratchet teeth and are indexed to succeeding teeth without imparting additional rotation to the auger. In the next cycle, the ratcheting action is repeated so that the auger rotates solely in one direction as the agitator oscillates in two directions.\nSince the pawls and teeth are disposed in a common horizontal plane and require initial pivoting of the pawls to engage the teeth, lost motion occurs therefrom. Accordingly, optimum efficiency of conversion of the oscillating agitator rotation to the unidirectional rotation of the auger is not achieved, which correspondingly decreases the efficiency of the washing operation."}
{"text": "The semiconductor integrated circuit (IC) industry has experienced rapid growth. Technological advances in IC materials and design have produced generations of ICs where each generation has smaller and more complex circuits than the previous generation. However, these advances have increased the complexity of processing and manufacturing ICs and, for these advances to be realized, similar developments in IC processing and manufacturing are needed. In the course of IC evolution, functional density (i.e., the number of interconnected devices per chip area) has generally increased while geometry size (i.e., the smallest component (or line) that can be created using a fabrication process) has decreased.\nTo achieve the shrinking geometry sizes, advanced lithography processes have been developed. During the lithography process, a general purpose (GP) nozzle is used to dispense a developer solution on a semiconductor wafer. Over time, the tip of the GP nozzle may have a buildup of water drop or water residue, which can adversely affect the intended operation of the GP nozzle. Consequently, semiconductor fabrication performance may be degraded.\nTherefore, while existing lithography processes have been generally adequate for their intended purposes, they are not entirely satisfactory in every aspect."}
{"text": "1. Field of the Invention\nThe present invention is broadly concerned with an improved hydraulic portioning assembly designed to facilitate precise portioning of fluid or semi-fluid products such as meat in order to deliver successive portions to the filling or packaging device. More particularly, it is concerned with such a portioning assembly making use of hydraulic control means comprising a double-acting piston and cylinder assembly coupled via valve means with a hydraulically powered product dispensing device such that stroke length of the hydraulic control means controls the portions dispensed.\n2. Description of the Prior Art\nIn the production of meat products such as hamburger or sausage, it is often desirable to deliver precisely metered or portioned quantities to a filling or packaging station. As can be appreciated, if underweight portions are dispensed, the ultimate consumer suffers. By the same token, if portions are overweight, the seller is at a disadvantage. Thus, it is very desirable in automated meat handling operations to provide a portioning assembly which gives precise, repeatable portions.\nBy the same token, however, modern day automated processing demands a given portioning assembly be adaptable for high speed operation and be extremely reliable in use. These requisites are of course inimical to precision portioning, and thus there has been a real need in the art for a simplified, rugged, yet highly precise portioning assembly useful in conjunction with hydraulic food pumps or similar devices."}
{"text": "The statements in this section merely provide background information related to the present disclosure and are not admitted as prior art.\nMoving Picture Experts Group (MPEG) and Video Coding Experts Group (VCEG) have developed video compression technology that is superior to the existing MPEG-4 Part 2 and H.263 standard. This new standard is called H.264/AVC (advanced video coding) and was jointly announced as MPEG-4 Part 10 AVC and ITU-T Recommendation H.264.\nIn H.264/AVC (hereinafter, simply referred to as H.264) standard, a residual signal is generated by performing an intra/inter prediction process in units of macroblocks having various types of sub-blocks, and encoding is performed by frequency transform, quantization, and entropy encoding in frequency transform units of 4×4 units with respect to the generated residual signal.\nVideo compression technologies are currently in development to meet the needs for ultra-high video resolution. The international standard organizations, MPEG and VCEG have jointly developed HEVC moving picture encoding standard in the name of JCT. In the ultra-high-resolution compression, experiments have proved that, in addition to 4×4 and 8×8 frequency units having been widely used in the existing video compression technologies, large frequency transform unit such as 16×16 is helpful to improve video compression efficiency. However, if such 16×16 unit is scanned and encoded at a time in the existing method, devising and implementing efficient frequency coefficient encoding is difficult."}
{"text": "1. Field of the Invention\nThe present invention relates to a tag structure for attachment to an article cast from molten material, to a method for labeling an article cast from a molten material and to an article cast from molten material including a tag having human or machine readable markings relative to the article, such as article identification.\nUpon completing the manufacture of an article cast from a molten material, there is a need to provide the subsequent user with a clear identification and other information relating to the article. For example, if the article is in its usable final form, product information such as price and weight may be useful. If the article is utilized in combination with various other items, a clear description of its individual function may be required. Further, production schedules may require inventory information to be displayed on the article. If the article is simply in an intermediate form to be subjected to further processing, information regarding the individual article's composition may be required for the final processing steps.\n2. Description of The Prior Art\nCertain methods are known for providing a display of information on articles in general. One method, perhaps the most common method, includes the application of a label coated with an adhesive backing which is adhered to the article. Unfortunately, many of these labels lack sufficient adherence to provide a reliable bond between the label and the article for extended periods of time.\nAnother method includes the attachment of a label in the form of a structural tag to the article by physically driving retaining projections which are integral with the tag into the article. An example of such a tag and process is disclosed in U.S. Pat. No. 3,673,717. Use of tags with these projections in such a process may subject the article to unwanted damage or at least create stresses within the article. Further, the hardness of the article to be labeled may preclude use of this method. This is especially likely in articles which are cast from molten material, such as aluminum, for example.\nAnother method, specific to cast molten material articles, includes the attachment of a label or tag before solidification of the molten material. An example of such a process is disclosed in U.S. Pat. No. 1,561,427. The process disclosed in this patent includes assembling individual characters on a strip to form a number or word, embedding the entire strip into a mold and casting the metal in the mold.\nRegarding this process, it should be noted that a number of individual characters must be arranged to create the marking. This may be time consuming and laborious, especially if much information is to be conveyed. Further, computer bar codes which are often used in many of today's products cannot be embedded in this manner. Also, repeated impression of the markings into the mold may eventually cause damage to the mold surface.\nStenciling and riveting are also known methods used with cast articles after the articles are made.\nIt can be seen that labeling an article can proceed during the making of the article or after the article is made, with the latter being the predominate mode."}
{"text": "As a boat is connected to a mooring buoy there are conditions where the boat can come into contact with the mooring buoy. This contact, depending on the conditions, can cause significant damage to the boat. For example, there may be damage caused by impact and abrasion between the boat and mooring buoy. It is also possible to encounter this same damage as one is attempting to dock at virtually any type of buoy.\nAccordingly, it is an object of the present invention to minimize the potential damage caused by contact between the vessel and the buoy.\nAnother object of the present invention is to provide a mooring buoy cover that prevents damage to the vessel, that can be applied easily to the buoy and that can be fabricated relatively inexpensively."}
{"text": "1. Field of the Invention:\nThis invention relates generally to a method and apparatus for determining the location of a point and more particularly to a method and apparatus for determining the coordinates of a contact point on a resistive type semianalog sensitive surface.\n2. Discussion of Background:\nThere are known, such as shown by French Pat. No. 2,520,498 sensitive surfaces of the type consisting of a support, generally of a glass plate, on which are placed a grid consisting of two terminals between which a certain number of emitting resistant strips are connected in parallel. Also placed on the support are conductive detecting strips placed between the resistant strips and connected separately to a selecting switch so that the voltage of each conductive strip is sampled indidually in successsion. The emitting resistance strips and the conductive detecting strips, which can also have some resistance, most often consist of a thin conductive enamel deposited on the glass plate. When a conductive object, for example an ordinary finger, comes to rest on this sensitive surface, it forms an electrical connection between at least one detecting strip and at least one emitting strip of the grid. By scanning the various strips with a selector, it is possible to digitally determine the coordinates, which are determined by the number of the detecting strip where a voltage is detected, and an analog determination of the other coordinate by the value of this voltage.\nPreferably, as described in said French patent application, said emitting strips are alternately polarized in opposite directions. For each direction of polarization, a certain number of measurements are made of the voltage of each detecting strip whose averages are computed separately, and the coordinates of the finger are deduced by comparing the averages obtained.\nAlthough such a process makes it possible to solve a great number of problems presented by the use of sensitive surfaces of this type, some difficulties in practical use still remain, such as are listed below.\nThe electrical contact resistance of the skin at the end of the operator's finger is highly variable and most often large (50 k.OMEGA. to 1M.OMEGA.) and would allow only the passage of a current of very small intensity between the grid and the detecting strip. Therefore, only a detector with very high impedance can be used.\nFor obvious reasons of user safety, only very low excitation voltages can be used. The user's body conducts direct voltages through the fingertip not only from contact with another point of the circuit, but also by static electricity, and especially alternating voltages by an antenna effect with ambient spurious signals. The combination of these two effects makes it difficult to obtain a favorable signal/noise ratio, and consequently, precise measurements.\nThe variability of the useful signal in relation to a constant background noise leads to extremely variable relative variations over the extent of the surface, particularly in the direction of the analog measurement nearing zero.\nFinally, scanning the various strips by the selector and taking various measurements can take considerable time, which is impractical for the user if standard filtering techniques are used.\nThe present invention aims at providing a method and apparatus that make it possible to overcome the above difficulties which are encountered in scanning of sensitive surfaces of this type, i.e., which meet all the requirements of safety, precision and speed."}
{"text": "1. Field of the Disclosure\nThe disclosure relates to a driving assistance apparatus for a vehicle, the driving assistance apparatus executing following control, which is control for causing a host vehicle to follow a preceding vehicle.\n2. Description of Related Art\nAs proposed in Japanese Patent Application Publication No. 2008-2456 (JP 2008-2456 A), a drive mode can be switched between a plurality of modes including a normal mode, an eco-mode, and a sport mode. The drive modes specify output characteristics of an engine relative to accelerator depression amounts. In the eco-mode, the engine output relative to the accelerator depression amount is limited in comparison with the normal mode in order to conserve fuel. In the sport mode, on the other hand, the engine output relative to the accelerator depression amount is set to be slightly larger than in the normal mode so as to improve a responsiveness by which the engine output responds to an accelerator operation. The drive mode is switched by a selection operation performed by a driver (i.e., in accordance with preferences of the driver).\nIncidentally, a driving assistance apparatus that causes a host vehicle to travel so as to follow a preceding vehicle travelling in front of the host vehicle may be employed in a vehicle to lighten a driving operation load of the driver. This control for causing the host vehicle to follow the preceding vehicle will be referred to as following control. During the following control, a target acceleration that enables the host vehicle to follow the preceding vehicle is calculated such that a vehicle-to-vehicle distance between the host vehicle and the preceding vehicle is maintained at a distance within a predetermined range, whereupon an engine or a brake apparatus is controlled on the basis of the target acceleration.\nIn a vehicle having both switchable drive modes and a following control function, however, the drive mode selected by the driver is not reflected during execution of the following control. In other words, the target acceleration calculated during the following control is calculated without taking the drive mode into account, and therefore the following control cannot be implemented in accordance with the preferences of the driver."}
{"text": "Electron emission devices are key components of many modern technological products. For example, focused \"beams\" of electrons produced by such devices are used in X-ray equipment, high-vacuum gauges, televisions, large-area stadium-type displays, and electron beam analytical devices such as scanning electron microscopes.\nStandard electron emission devices operate by drawing electrons from a cathode formed from a material that readily releases electrons when stimulated in a known manner. Typically, electrons are drawn from the cathode by the application of either a thermionic stimulus or an electric field to the cathode. Devices operating through application of an electric field are said to operate by field emission. Cathodes used in field emission devices are accordingly known as field emission cathodes, and are considered \"cold\" cathodes, as they do not require the use of a heat source to operate.\nField emission offers several advantages over thermionic stimulus in many electron emission applications. A field emission device (which creates an electric field) will typically require less power than a thermionic device (which creates a heat source) to produce the same emission current, respectively. Field emission sources are typically on the order of 1000 times brighter than comparable thermionic sources. The added brightness can be highly advantageous in lighting applications, such as stadium displays, or in applications which require the use of electron beams operating at intense focus, such as microscopes.\nFurther, the heat sources used in thermionic electron emission devices eventually damage them, leading to relatively quick \"burnout.\" In applications requiring the use of many electron emission devices, such as in large area collective usage television screens, use of thermionic emission devices is very expensive because of the need to replace frequently devices suffering from rapid burnout.\nAdditionally, thermionic electron emission devices are not feasible for some applications. Thermionic devices are temperature dependent, and thus cannot be used in applications operating in extreme temperatures or where the ambient temperature conditions vary substantially over time. For example, thermionic devices will not work properly in motors or engines where temperature conditions may swing from 70.degree. Fahrenheit to -60.degree. Fahrenheit within a few minutes. In contrast, field emission devices, which operate relatively independently from temperature conditions, can be used in such applications. Thermionic devices are also inappropriate for use where the heat used to draw the electron beam may damage the environment within which electron emission is to occur. For example, in X-ray applications focused near a human body, thermionic emission of electrons is undesirable as the heat source applied could cause pain or damage to the subject. Field emission devices avoid these concerns as they apply and generate relatively little heat.\nAmong various materials known to be suitable for the construction of field emission cathodes, carbon-based materials have proven to be capable of producing significant emission currents over a long lifetime in relatively low-vacuum environments (10.sup.-7 Torr or less). Cathodes utilizing diamond films, bulk carbon, and graphite have been developed, but have required the application of substantial voltages to the cathode before generating significant electron emission. Other cathodes having regular, defined surface structures created from carbon materials include cathodes constructed from individual carbon fibers bundled together, cathodes machined from carbon rods, and matrix cathodes with carbon surfaces formed by photolithography and thermochemical etching procedures. While these cathodes can produce high current density upon application of low voltages, they are expensive to produce, as they require sophisticated fabrication procedures and/or manual assembly in their production.\nIt is an object of the current invention to provide an efficient and durable field emission cathode which may be manufactured simply and inexpensively.\nAnother object of the current invention is to provide a field emission cathode comprising an emissive member formed of a porous carbon foam material having an emissive surface which defines a multiplicity of emissive edges.\nOther objects and advantages of the current invention will become apparent when the field emission cathode of the present invention is considered in conjunction with the accompanying drawings, specification, and claims."}
{"text": "1. Field\nThe invention relates to a method of flying an unmanned aerial vehicle (UAV) in response to emergency conditions.\n2. Background\nUAVs are typically controlled by a remote operator via wireless communication. It is often the case that when a component of the UAV fails, or some functionality is lost, the unmanned aerial vehicle is unable to carry out a suitable change to its flight plan.\nFurthermore, nothing in the prior art considers the possibility of multiple failures occurring.\nTherefore, it would be desirable to have a method and apparatus that takes into account at least some of the issues discussed above, as well as other possible issues."}
{"text": "This invention relates to the deposition of thin films and is more particularly concerned with a reactive sputtering process wherein atoms of a target material are freed from a conductive target to form a coating which is deposited on a substrate. This process can be employed, for example, in a reactive sputtering process, for creating dielectric insulating layers on electrical parts, or wear-resistant layers on mechanical parts.\nThe invention is more specifically directed to an improved dc sputtering process in which the dielectric coating material that becomes lodged on the conductive target is removed, thus avoiding a major cause of arcing.\nReactive sputtering, and problems related to dc sputtering, are described in great detail in my earlier U.S. patent application Ser. No. 08/261,988, filed Jun. 17, 1994, and having a common assignee herewith.\nSputtering is a vacuum deposition process in which a sputtering target is bombarded with ions, typically an ionized noble gas, and the atoms of the target material are mechanically freed by momentum transfer. The target material then coats a nearby substrate.\nIn a reactive sputtering process a reactive gas is introduced into the deposition chamber, and the freed target material reacts with the reactive gas to form a coating material. For example, the target material can be aluminum, and oxygen can be introduced as the reactive gas to produce a coating of aluminum oxide. A carbonaceous gas, e.g. acetylene, can be used as the reactive gas to produce carbide coatings such as SiC, WC, etc., and N.sub.2 can be introduced to produce a nitride coating such as TiN. In any event, the conductive target atoms and the reactive gas react in plasma in the chamber to produce the compound that serves as a coating. In a typical example, aluminum atoms freed from an aluminum target enter a plasma of argon and oxygen to produce a deposition of aluminum oxide.\nDC sputtering is a random process and the insulating coating material is deposited on all available surfaces. This means that not only does the insulating material coat the article in question, but it also coats other surfaces in the chamber as well, including the target. Thus, in a reactive sputtering process for depositing aluminum oxide, molecules of Al.sub.2 O.sub.3 land on the surface of the aluminum target. This deposition of an insulator on the target causes severe problems, including a reduction of sputtering rate and a propensity to induce arcing.\nContamination of the target can also result, even in conventional dc sputtering, due to atmospheric gases, water droplets, inclusions, and other contaminants. Each of these can be a source for arcing, and their presence will also reduce the deposition rate over time because of a reduced active sputtering area on the target. Accordingly, these problems necessitate a frequent cleaning of the target surface.\nThis problem has been known to exist for some time, but its causes have not been completely appreciated. Procedures to deal with these problems such as arc control in active sputtering have not been completely satisfactory.\nA standard approach involves sensing the presence of arcing, and then interrupting current flow. This will control the arcing, but does nothing about the insulating coating that continues to cover the target.\nOne early attempt to deal with arcing in a blind fashion involves periodically interrupting the current flow between the dc power supply and the plasma generator in which the sputtering occurs. Here, shutting off the dc power serves to extinguish early arcing. This means that unipolar pulses of power at a fixed duty cycle are fed to the target. This does have the beneficial effect of permitting charge to build up only partially across the dielectric deposition on the target, so that arcing is less likely to occur, and also can lead to a small amount of resputtering of the deposition. However, this system, while slowing down the rate of insulating deposit on the target, does not reverse the deposition.\nAnother system previously proposed is called a low energy small package arc repression circuit. In that system, an electronic switch cycles at a rate of about 2 KHz to cut off the current to the target. This in effect reverses the voltage on the target to several volts positive, and draws some electrons from the plasma to the from surface of the insulative deposition. This neutralizes the anions on the front surface of the deposition, to discharge the voltage buildup across the layer, thereby greatly reducing the occurrence of dielectric breakdown and arcing. Also, discharge of the front surface of the insulative layer lowers the surface potential to approximately that of the target. Discharging of the dielectric deposition also permits the argon ions in the plasma to collide with the insulating dielectric material. This does result in some resputtering of the molecules of the deposited material, thus slowing the rate of deposit on the target. However, this approach does not resputter the molecules of the deposited compound as effectively as the atoms of the target material, and so this approach has not been totally effective in removing deposits from the target during reactive sputtering processes.\nDifferent materials require different voltages to be applied to the targets to effect sputtering. For example, because a gold atom is a much heavier atom than an aluminum atom, it requires a much more energetic ion to free it from the target. Typically, in a process that employs an aluminum target, the applied voltage needed is about -450 volts, while in a similar process that employs a gold target the applied voltage has to be about -700 volts.\nConsidering that an aluminum oxide (Al.sub.2 O.sub.3) molecule is significantly heavier than an aluminum atom, one can understand that a higher potential would be required to energize the argon ion enough to resputter the coating. This is, of course, true for other materials as well.\nAnother approach to solve this problem involves a pair of sputtering targets, with one serving as cathode and the other as the anode. The applied electrical voltage is periodically reversed so that the sputtering occurs first from the one target and then from the other. This process also reverses the charge on the deposited insulating material as well, which reduces the possibility of arcing and also resputters some of the insulating material on the targets. However, this arrangement, requiring plural targets, can be cumbersome and expensive to employ. Moreover, this arrangement suffers from a low deposition rate due to the amount of time needed for crossover.\nThese previous solutions, which employ unipolar pulsing or alternately cycled targets, have been somewhat effective in reducing voltage stress on insulating films redeposited on the targets, but have not been completely effective in removing the redeposit or in preventing it. None of these approaches sputters off the insulator from the beginning before it has a chance to accumulate, and none of these techniques has been effective in eliminating or halting the redeposition of insulating film on the target, while maintaining a high deposition rate.\nIn the process disclosed in my earlier patent application Ser. No. 08/261,988, reactive dc sputtering is carried out in a plasma chamber by applying an electric potential to a conductive target so that target material is sputtered from the target and is reacted with a reactive gas that is introduced into the chamber. A dc voltage of a suitable level, e.g., -500 volts, is applied from a power supply to the sputtering target, which serves as cathode, with a conductive surface in the plasma chamber being held at ground to serve as anode. A noble gas, e.g., argon, present in the chamber is ionized and creates a plasma, e.g., argon anions and free electrons. The electrons are drawn off in the anode, and the positive argon ions are accelerated towards the cathode, that is, towards the conductive target. The argon ions knock atoms of target material free from the target by momentum transfer. The argon ions pick up electrons from the negatively charged target, and migrate back to the plasma. The freed target atoms enter the plasma and react with a reactive gas that has been introduced into the chamber. The reactive gas can be, for example, oxygen, nitrogen, borane, acetylene, ammonia, silane, arsene, or various other gases. The reaction product is deposited on a substrate positioned adjacent the plasma. The substrate can be a masked semiconductor wafer on which a compound such as a Al.sub.2 O.sub.3, SiO.sub.2, or another insulator or dielectric is to be deposited. In some processes, the substrate can be a drill bit, wear plate, valve spindle, or other mechanical part on which a wear resistant coating, such as WC or TiN, is deposited.\nNon-reactive sputtering can be used to deposit a conductive layer onto the master in a compact disk fabrication process, to permit subsequent electroplating.\nAs mentioned before, the reaction product of the reactive sputtering process is deposited randomly, and not only coats the workpiece substrate, but also coats other surfaces including walls of the chamber and the sputtering target. Accumulation of the insulating coating can induce arcing, and also reduces the available area of the sputtering target, thereby reducing the sputtering rate over time.\nIn the process described in my earlier patent application Ser. No. 08/261,988, the dc power that is applied to the target is interrupted periodically by reverse-bias pulses of a dc voltage level that is positive relative to the anode. Preferably, the reverse bias pulses are at a level of 50 to 300 volts above ground potential, and these are applied at a pulse frequency of 40 KHz to 250 KHz with a pulse width of 1 .mu.sec to 3 .mu.sec. This results in low duty cycle pulses (about 10% or less). The reverse bias creates a reversal of the charge across the insulating material. These accumulations behave as a capacitor, with the conductive target being one plate and the conductive plasma being the other plate. The reverse voltage is applied long enough (e.g. 2 .mu.sec) for the polarity of the capacitive charge to be reversed, up to -300 volts, on the plasma side of the deposition.\nThen when the normal or negative sputtering voltage is again applied, the argon ions in the plasma are accelerated preferably towards the reverse-charged dielectric material. These ions are accelerated to an increased energy due to the additional potential difference. As a result, the molecules of the deposition are re-sputtered off the target. This process also permits the active sputtering surface area of the target to remain as large as possible.\nThis preferential sputtering process also cleans out other insulative contaminants from the surface of the target whether used for reactive sputtering or conventional sputtering.\nAlthough the aforedescribed preferential sputtering process has largely achieved favorable results in the field of pulsed dc power, additional problems can arise that this technique does not resolve. For example, particulates and arc damage to the surface of the target can give rise to mechanical arcs, that is, persistent arcing problems that do not clear themselves immediately when the reverse bias is applied. Also, an extreme voltage condition can exist, for example if there is too little gas present to support a plasma. In most cases an arcing or an overvoltage condition will clear itself in a single cycle of the applied pulsed voltage. However, where the reverse bias voltage does not clear the problem, some additional steps need to be taken to prevent damage either to the workpiece, the plasma cell, or the power supply. Unfortunately, no present system is capable of operating in conjunction with the reverse bias pulsed power supply and of detecting and reacting to either a persistent arc (sticky arc or multistrike arc) or an overvoltage condition."}
{"text": "This invention relates to a liquid level reduction alarm device for a liquid reservoir vessel kept in liquid communication with a braking master cylinder, preferably in automotive use and adapted for delivery of alarm signal to the vehicle driver, if the liquid level in the reservoir vessel should reduce to an excessive degree, so as to warn him to replenish the brake oil.\nIt is a main object of the present invention to provide a substantially shock-proof electric and built-in alarm system highly effective in the aforementioned field of engineering.\nA further object is to provide an effective alarm system of the above kind which is operable magnetically and automatically responsive to the liquid level contained in the reservoir vessel.\nA still further object is to provide an alarm system of the above kind which is highly effective to avoid misoperation caused by outside vehicle shocks.\nThese and further objects, features and advantages of the present invention will become more apparent from the following detailed description of the invention by reference to the accompanying drawing illustrative of a preferred embodiment of the present invention."}
{"text": "Large organizations need to do a great deal of electronic record keeping, utilizing large and complex systems, requiring extensive hardware and software infrastructure. Example systems include a university system for tracking applications and enrolled student histories, an insurance company system for tracking clients and insurance claims, or a police department system for tracking information relevant to active cases and to individual criminal histories. Each of these systems may be used by many people simultaneously, possibly spread across a large physical area, and need to track a large number of documents across a long period of time. Each system additionally requires specific software to serve the needs of the organization effectively. When a large organization commissions a new system of this type, typically a group of application design, system administrator, and system architecture experts is brought in to design the software and the infrastructure (e.g., the hardware on which the software runs), and then to scale the system into production use. Designing the infrastructure is a complex task, and depends on the needs of the organization in many ways. Each resulting system is different, causing it to take a great deal of time to thoroughly build and test new systems. When problems occur and software support becomes involved, it may take a significant effort even to communicate the relevant details of the infrastructure setup so troubleshooting can begin."}
{"text": "The present invention is directed to reflectors used in the radiant section of a fired heater, and more particularly to radiant reflectors provided on a refractory wall centered in the spacing between the radiant tubes.\nCombustion equipment is generally operated in chemical plants, petrochemical plants and refineries. The equipment may include industrial heaters, furnaces or plant boilers. This equipment is generally designed with bare or smooth-walled tubes, or with partially studded tubes as disclosed in my earlier U.S. Pat. No. 6,364,658, which is hereby incorporated herein by reference in its entirety. Use of tubes in radiant sections usually exposes the front half of the tube to direct flame radiation, while limiting the exposure of the rear half or dark side of the tube to reflected radiation.\nThe heat flux distribution around the circumference of a conventionally fired tube at a conventional spacing of two tube diameters is depicted in FIG. 1. A flame or radiating plane is on one side of the tube and a refractory wall is on the other. The front half of the tube surface faces the flame (point A) and receives a higher heat flux as compared to the rear half facing the refractory wall (point B). Point A receives heat flux only from direct flame radiation, while point B, facing the refractory wall, receives only reflected radiation coming from the refractory wall. Points between point A and point B receive varying amounts of both direct and reflected radiation, depending upon their location along the tube.\nThe standard distance between tubes is two tube diameters from center-to-center, and 1.5 diameters from the center of the tubes to the refractory wall, for most operations in the chemical and petrochemical industries, as shown in FIG. 2. The heat flux distribution in FIG. 1 is based on this configuration. For the purposes of an illustration using fluxes typical in a conventional fired heater, where the highest heat flux at point A is 18000 Btu/hr-ft2, the diametrically opposed counterpart (point B) receives only 6000 Btu/hr-ft2. The rear half of the tube transfers only 24% of the total heat absorbed by the tube; this includes both the direct and reflected radiation, as seen in FIG. 3. The average flux for the tube amounts to 10,000 Btu/hr-ft2.\nMore than 85% of the heaters in the industry have such a large flux differential between the front and the rear side of the tube, as this illustration depicts. A significant compromise is made on the overall heat-receiving capacity of the tube in order to keep the flame-front side (point A) within safe working temperatures.\nTo make the heat flux distribution in the tube more uniform, one approach of the furnace designers has been to increase the center-to-center tube spacing requirements from 2 to 3 tube diameters. This design increases the flux at point B of the tube from 6,000 Btu/hr-ft2 to 9,000 Btu/hr-ft2 as shown in FIGS. 4A and 4B. The increased spacing has the beneficial result of increasing the heat-receiving capacity of the rear half of the tube for the 3D-spaced tubes, while heat flux distribution on the front half of the tube is generally the same as for the 2D-spaced tubes. This results in an increase of the average heat flux to 12,000 Btu/hr-ft2 for the entire tube. However, the drawback of this solution is apparent. With an increase in tube spacing there is a corresponding increase in the size of the heater. This increases the cost and space requirements for the heater.\nAnother prior art approach improves the heat flux distribution by placing radiating flames on opposing sides of the tubes in a so-called xe2x80x9cdouble-firedxe2x80x9d design. A comparison is shown between one radiating flame (A) and two radiating flames (B) in FIGS. 5A and 5B, respectively. This design is commonly used in chemical processes that mandate a more uniform heat flux distribution, such as, for example, in delayed cokers, high-pressure hydrotreaters, ethylene furnaces, and the like. In a double-fired system, the front (point A) and rear (point B) portions of the tube have the same heat flux rate due to direct flame radiation, and the points at the margins between the front and rear receive relatively less direct flame radiation. The corresponding distribution of the heat flux, for the illustrative example, is 18,000 Btu/hr-ft2 for the front and the rear locations, 13,500 Btu/hr-ft2 at the margins between the front and rear faces, i.e. the middle area of the tube (point M at the 90 and 270 degree positions), resulting in an average flux of 15,000 Btu/hr-ft2. The double-fired design brings with it the disadvantage that the heater has to be much larger, as much as twice the size as a single-fired unit, and correspondingly more expensive.\nThe present state of technology for heaters with a standard spacing of 2 tube-diameters will have a relative flux ratio of 1 to 1.8 between the average flux and the maximum flux, whereas a heater with a 3 tube-diameter spacing will have a relative flux ratio of 1 to 1.5, as shown in API Standard 530, Calculation of Heater-Tube Thickness in Petroleum Refineries, American Petroleum Institute (1988), Figure C-1 Ratio of Maximum Local to Average Heat Flux Curves, page 103.\nThe 3 tube-diameter design is less common in the industry and the vessel must be significantly larger than a 2 tube-diameter design. The average to maximum flux ratio of the double-fired tubes is significantly lower at 1 to 1.2, but is a more costly alternative of the three designs for an industrial plant.\nA recent improvement in the flux distribution as described in my U.S. Pat. No. 6,364,658, involves the placement of extended surfaces such as studs or fins on the dark side of the tubes in a single-fired arrangement. This improves the heat transfer to the dark side of the tubes primarily by increasing the convection heat transfer. Still, in the standard tube arrangement with smooth walls, it is well known that 65.8% of the radiant heat from the flame is absorbed by the tubes, primarily the front half of the tubes facing the flame, and 34.2% goes through the spaces between the tubes to the refractory wall. The same percentages apply to the reflected radiation from the refractory onto the dark side of the tubes, i.e. 65.8% of the 34.2% is re-radiated to the rear half of the tubes, or 22.5%. In other words, 88.3% is absorbed by the tubes, front and back, and the balance of 11.7% is radiated back to the flame through the spaces between the tubes. It would be very desirable if a significant portion of this 11.7% could be directed onto the tubes instead of the flame. There thus remains a need for making the flux distribution even more uniform and/or for increasing the rate of heat absorption by the tubes.\nThe present invention utilizes radiation reflectors positioned on the refractory wall of a furnace, preferably in the spaces between the radiant tubes. The radiation reflectors provide surfaces which are angled, with respect to generally flat or curvilinear refractory surfaces behind the tubes, to reduce the radiation that is reflected between the tubes and increase the radiation reflected onto the dark side of the tubes. The use of the radiation reflectors thus increases the radiant flux delivered to the dark side of the tubes, increasing heat absorption and decreasing the ratio of the maximum to average flux. The radiation reflectors can also enhance convection heat transfer to the dark side of the tubes by increasing the velocity of the flue gases between the tubes and the refractory wall, thereby increasing the convection heat transfer.\nIn one aspect, the present invention provides radiation reflectors for use in a fired furnace comprising a plurality of parallel tubes arranged in a row between a flame on a radiant side and a generally flat or curvilinear refractory surface on a dark side. The radiation reflectors have a longitudinal base for abutment against the refractory surface. The base has opposite edges at either side thereof. A longitudinal cusp is opposite the base, and longitudinal reflective surfaces extend from each edge of the base to the cusp. The reflective surfaces have concavity in a plane transverse to a longitudinal axis, preferably parabolic sections in the transverse plane. An anchoring pin can extend transversely through each radiation reflector from the cusp into a subjacent structure.\nIn another aspect, the invention provides a fired furnace for heating petroleum, petrochemicals or chemicals. The furnace has a plurality of parallel tubes each disposed in a row between a flame on a radiant side thereof and a refractory surface on a dark side thereof. There are spaces between adjacent tubes. Radiation reflectors are positioned on the refractory surface opposite the spaces to reflect incident radiation from the flame away from the spaces and onto the dark side of the tubes. A central longitudinal bore is provided through each tube for the passage therethrough of a fluid to be heated. The row of tubes can be straight or circular. The radiation reflectors can be disposed longitudinally on either side of a flat surface of the refractory surface opposite a tube.\nIn a further aspect, the invention provides an improvement in a fired furnace. The furnace includes a plurality of parallel tubes disposed between a flame and a refractory wall. Adjacent tubes define a space between the tubes, and each tube includes a central longitudinal bore for the passage therethrough of a fluid to be heated and an outside diameter having a radiant side for exposure to radiation from the flame and a dark side essentially free of direct exposure to the flame. The improvement comprises positioning the radiation reflectors described above on the refractory wall opposite each space. Preferably, the reflective surfaces are parabolic sections in the transverse plane focused on the dark side of the adjacent tubes.\nIn a still further aspect of the invention, there is provided a method for improving the heat transfer in a fired furnace comprising a plurality of parallel tubes disposed between a flame and a refractory wall. Adjacent tubes define spaces between the tubes. The refractory wall comprises a generally flat or curvilinear surface opposite the tubes and spaces. The method includes the step of installing the radiation reflectors described above on the refractory wall opposite the spaces. The installation can include pinning the radiation reflectors with a pin extending from the cusp into the refractory wall. The radiation reflectors are preferably focused to reflect incident radiation from the flame onto the adjacent tubes on either side of a respective space. The tubes can have extended surfaces at least on the dark side. Where the tubes have smooth outside walls, the method can also include removing the smooth-walled tubes from the furnace and replacing them with tubes that have extended surfaces on a dark side opposite the refractory."}
{"text": "Modern communication networks may be used to deliver media content, such as movies, television shows, or video-on-demand programs, to subscribers. A service provider may also provide the subscriber with an interactive display. One problem with presenting an interactive display relates to determining when to display the interactive display. For example, when the interactive display is displayed, the interactive display should not interrupt the media content being delivered to the subscriber or prevent the subscriber from viewing an entire media content segment. Further, when the subscriber is interacting with the interactive display, the interaction should not prevent the subscriber from viewing subsequent media content. Thus, there is a need for an improved system and method for displaying media content and an interactive display."}
{"text": "The most varied membrane materials and adsorbents are known for separating mixtures of materials and are continually being improved with respect to their technical suitability and economy. In the treatment of effluent waters, it is desirable for instance by developing asymmetrical membranes to achieve the throughput loadings necessary for technical application. On the other hand, other problems arise in other areas of use. For example, known membrane materials, such as cellulose acetate membranes, are not very temperature or pressure resistant and swell rapidly in organic solvents. The low temperature, pressure and solvent resistance has the result that the pore sizes continually alter during technical application and this can lead to non-reproducible results as well as short periods of use for the membranes.\nAlso, adsorbents for purifying fluid or gaseous media are known in enormous numbers, for example, active charcoal and oxides with high specific surfaces, such as aluminium oxide. With these known adsorbents, adsorption normally takes place non-specifically and not very selectively. Moreover, the adsorption process is often only poorly reversible at high adsorption grades, that is, a high consumption of energy is necessary for desorption.\nKnown adsorbents and membranes with surfaces which have been chemically modified also have the disadvantage that, as a rule, they have only a modified mono-layer at the surface, so that they are extraordinarily sensitive to mechanical and chemical effects, such as hydrolytic action.\nThe inventors have already developed certain silicic acid heteropolycondensates, which are particularly suitable for the preparation of temperature, pressure and solvent-resistant membranes for the most varied separation problems. In the granular or similar state, these polycondensates are usable as adsorbents, even for highly specific separations, which either cannot be carried out or only with unsatisfactory results with known adsorption agents.\nThese porous silicic acid heteropolycondensates usable as membranes and adsorbents are made by:\n(i) forming a reaction mixture comprising: PA0 (ii) subjecting the reaction mixture to hydrolysis and condensation in the presence of at least the amount of water stoichiometrically required for hydrolysis, optionally a solvent and in the presence of a condensation catalyst, wherein the amounts of the components (a) to (d) are selected so that the resultant silicic acid heteropolycondensate, calculated as oxide units, contains, by weight, 35% to 90% of component (a), 10% to 50% of component (b), 0% to 15% of component (c) and 0% to 40% of component (d).\n(a) at least one hydrolysable silicic acid derivative of the general formula I: EQU SiR.sub.4 (I) PA1 wherein each R is selected from the group consisting of hydrogen, halogen, alkoxy and --NR'.sub.2 radicals, wherein each R' in the last-mentioned radical independently represents a hydrogen atom or an alkyl group, with the proviso that not all the R radicals are hydrogen atoms, and PA1 (b) at least one substituted silane of the general formula II: EQU SiR.sub.n R\".sub.(4-n) (II) PA1 wherein R has the meaning defined above, each R\" is selected from the group consisting of alkyl, alkenyl, aryl and aralkyl radicals and n is an integer from 1 to 3, and PA1 (c) optionally at least one functional silane of the general formula III: EQU SiR.sub.n (R'\"Y).sub.(4-n) (III) PA1 wherein R and n have the meanings defined above, each R'\" is selected from the group consisting of alkylene, phenylene, alkylphenylene and alkylenephenylene radicals and Y is selected from the group consisting of halogen, amino, anilino, aldehyde, keto, carboxy, hydroxy, mercapto, cyano, hydroxyphenyl, diazo, carboxylic alkyl ester, sulphonic acid (--SO.sub.3 H) and phosphoric acid (--PO.sub.3 H.sub.2) radicals, and/or PA1 (d) optionally at least one oxide component selected from the group consisting of non-volatile oxides soluble in the reaction medium and compounds capable of forming non-volatile oxides soluble in the reaction medium, the oxide component being an oxide of an element selected from Groups Ia to Va, Group IVb and Group Vb of the Periodic System, and\nIn this method, the condensation catalyst is used in an amount of not more than 3% by weight, based on the weight of the reaction mixture."}
{"text": "Real power is the capacity of a circuit for performing work in a particular time. Apparent power is the product of the current and voltage of the circuit. The apparent power may be greater than the real power due to energy stored in the load and returned to the source, or due to a non-linear load that distorts the wave shape of the current drawn from the source.\nThe power factor of an AC electric power system may be defined as the ratio of the real power flowing to the load to the apparent power (a number between 0 and 1).\nIn an electric power system, a load with a low power factor draws more current than a load with a high power factor, for the same amount of useful power transferred. The higher currents may increase the energy lost in the distribution system, and may require larger wires and other equipment. Because of the costs of larger equipment and wasted energy, electrical utilities may charge a higher cost to customers with a low power factor.\nIn a purely resistive AC circuit, voltage and current waveforms are in phase, changing polarity at the same instant in each cycle. Where reactive loads are present, such as with capacitors or inductors, energy storage in the loads results in a time difference (phase) between the current and voltage waveforms. This stored energy returns to the source and is not available to do work at the load. Thus, a circuit with a low power factor will have higher currents to transfer a given quantity of real power compared to a circuit with a high power factor.\nAC power flow has the three components: real power (P) measured in watts (W); apparent power (S) measured in volt-amperes (VA); and reactive power (Q) measured in reactive volt-amperes (VAr). Power factor may thus be defined asP/S  (1)\nIn the case of a perfectly sinusoidal waveform, P, Q and S can be expressed as vectors that form a vector triangle such thatS2=P2+Q2  (2)\nIf θ is the phase angle between the current and voltage, then the power factor is equal to |cos θ|, andP=S*|cos θ|  (3)\nWhen power factor is equal to 0, the energy flow is entirely reactive, and stored energy in the load returns to the source on each cycle. When the power factor is equal to 1, all the energy supplied by the source is consumed by the load. Power factors may be stated as “leading” or “lagging” to indicate the sign of the phase angle.\nIf a purely resistive load is connected to a power supply, current and voltage will change polarity in phase, the power factor will be unity, and the electrical energy will flow in a single direction across the network in each cycle. Inductive loads such as transformers and motors consume power with the current waveform lagging the voltage. Capacitive loads such as capacitor banks or buried cables cause reactive power flow with the current waveform leading the voltage. Both types of loads will absorb energy during part of the AC cycle, which is stored in the device's magnetic or electric field, only to return this energy back to the source during the rest of the cycle. For example, to get 1 kW of real power, if the power factor is unity, 1 kVA of apparent power needs to be transferred (1 kW÷1=1 kVA). At low values of power factor, however, more apparent power needs to be transferred to get the same real power. To get 1 kW of real power at 0.2 power factor, 5 kVA of apparent power needs to be transferred (1 kW÷0.2=5 kVA)."}
{"text": "Radio frequency RF) receivers are used in a wide variety of applications such as television receivers, cellular telephones, pagers, global positioning system (UPS) receivers, cable moderns, cordless phones, satellite radio receivers, and the like. As used herein, a “radio frequency” signal means an electrical signal conveying useful information and having a frequency from about 3 kilohertz (kHz) to hundreds of gigahertz (GHz), regardless of the medium through which such signal is conveyed. Thus an RF signal may be transmitted through air, free space, coaxial cable, fiber optic cable, etc. One common type of RF receiver is the so-called superheterodyne receiver. A superheterodyne receiver mixes the desired data-carrying signal with the output of tunable oscillator to produce an output at a fixed intermediate frequency (IF). The fixed IF signal can then be conveniently filtered and converted down to baseband for further processing. Thus a superheterodyne receiver requires two mixing steps.\nModern integrated circuit technology has allowed many of the circuits used in RF receivers to be combined on-chip and thus to substantially reduce the cost of the RF receiver. However this level of integration creates other problems. For example, signals from one part of the chip may be electrically or magnetically coupled to circuits in another part of the chip. These unwanted signal couplings can distort the desired signal and create artifacts that can be perceived by the viewer or listener. Traditionally, integrated circuit designers have used layout strategies to reduce coupling between circuits, such as physical separation, the addition of ground rings, a reduction in the length of conductors, etc. However these techniques, while still useful, are unable to completely eliminate the deleterious effects of electrically or magnetically coupled energy.\nThe use of the same reference symbols in different drawings indicates similar or identical items."}
{"text": "In a conventional Multimedia Over Coaxial Alliance (MoCA) network, data packets are transmitted over a coaxial communication channel having a bandwidth of 50 MHz. FIG. 2 illustrates a conventional transmission processing channel 200 in a MoCA node. As shown in FIG. 2, data is first modulated at a modulator 202. Modulator 202 typically modulates the data using a high order modulation technique such as 16-bit quadrature amplitude modulation (16-QAM), 32-QAM, 64-QAM, 128-QAM, or 256-QAM. The modulated data is then power loaded and shaped at shaper 204 such that the signal spectrum exhibits an x/(sin x) frequency characteristic. Power loading includes adding power to particular subcarriers to compensate for the power roll off that occurs prior to transmission at the digital-to-analog converter (DAC) 212. After power loading and shaping, the data is transformed into the time domain using an Inverse Fast Fourier Transform (IFFT) 206. The output of the IFFT has 10-bit per sample precision which is input to a digital-to-analog converter (DAC) 212 to convert the data from a digital signal to an analog signal so it may be transmitted over the communication channel.\nEach data packet in a conventional MoCA network has a payload. A preamble is typically appended to the data packet. The preambles are used to calibrate the receiver to receive the following data with the appropriate gain and frequency compensation. The preambles are also used by the receiver to determine when the first OFDM symbol of the payload will be received at the receiver. FIG. 1 illustrates the five unique preambles of four different lengths that are utilized in MoCA networks. The preamble lengths are varied to maximize the available bandwidth for data transmission. As shown in FIG. 1, the two most robust preambles are the Beacon preamble and the MAP preamble. The Beacon and MAP preambles are the most robust and least efficient preambles because they are used for network coordination and contain the most critical information. The second most robust preamble type is the Admission Probes preamble. This preamble is used by a transmitter to accommodate a receiver that has little-to-no a-priori information about the link over which the transmitter will transmit data. The second most efficient preamble is the Broadcast/Data preamble. The Broadcast/Data preamble is used to transfer data from a transmitter to one or more receivers and assumes the receiving node has some a-priori information about the link. The most efficient, and thus shortest, preamble is the High-Throughput Unicast Data preamble. This preamble is appended to the beginning of a data packet when the receiver has a substantial amount of a-priori information about the communication channel.\nAs shown in FIG. 1, each of these five preambles is formed from five unique components. There are four long components, L1, L2, L3, and L4, and one short component, SS. Each of the long components is a 64-bit binary phase shift key (BPSK) modulated time domain bit sequence, and the short component is a 30-bit BPSK time domain bit sequence. The short sequence is the first component of a generic preamble and may be used by a receiving node to adjust the automatic gain control. Two of the long sequences, L1 and L2, may be used to finalize the automatic gain control adjustment, to perform burst detection, and to estimate the frequency offset between the transmitting node and the receiving node. The other two long sequences, L3 and L4, are used to identify link control packets to perform network coordination, admission, and maintenance. These components are generated in the time domain by modulating the following binary sequences:\n                    ⁢          S      =              {                  0          ,          1          ,          1          ,          1          ,          0          ,          0          ,          0          ,          0          ,          1          ,          0          ,          0          ,          0          ,          0          ,          1          ,          1          ,          1          ,          1          ,          1          ,          0          ,          0          ,          0          ,          1          ,          0          ,          0          ,          1          ,          1          ,          1          ,          1          ,          1          ,          0                }                        L      ⁢                          ⁢      1        =          {                                                  0              ,              0              ,              0              ,              0              ,              0              ,              0              ,              1              ,              0              ,              1              ,              0              ,              1              ,              0              ,              0              ,              1              ,              1              ,              0              ,              0              ,              1              ,              0              ,              0              ,              0              ,              1              ,              0              ,              0              ,              1              ,              0              ,              1              ,              1              ,              0              ,              1              ,              1              ,              0              ,                                                                          0              ,              0              ,              1              ,              1              ,              1              ,              0              ,              1              ,              0              ,              0              ,              0              ,              0              ,              1              ,              1              ,              0              ,              1              ,              0              ,              1              ,              1              ,              1              ,              0              ,              0              ,              1              ,              1              ,              1              ,              1              ,              0              ,              1              ,              1              ,              1              ,              1              ,              1              ,              0                                          }                  L      ⁢                          ⁢      2        =          {                                                  1              ,              0              ,              1              ,              1              ,              1              ,              0              ,              0              ,              1              ,              1              ,              1              ,              1              ,              0              ,              1              ,              1              ,              1              ,              1              ,              1              ,              0              ,              0              ,              0              ,              0              ,              0              ,              0              ,              1              ,              0              ,              1              ,              0              ,              1              ,              0              ,              0              ,              1              ,              1              ,                                                                          0              ,              0              ,              1              ,              0              ,              0              ,              0              ,              1              ,              0              ,              0              ,              1              ,              0              ,              1              ,              1              ,              0              ,              1              ,              1              ,              0              ,              0              ,              0              ,              1              ,              1              ,              1              ,              0              ,              1              ,              0              ,              0              ,              0              ,              0              ,              1              ,              1              ,              0              ,              1                                          }                  L      ⁢                          ⁢      3        =          {                                                  1              ,              1              ,              0              ,              1              ,              0              ,              1              ,              0              ,              1              ,              1              ,              1              ,              0              ,              0              ,              1              ,              1              ,              0              ,              0              ,              0              ,              0              ,              1              ,              0              ,              0              ,              0              ,              1              ,              0              ,              1              ,              0              ,              0              ,              1              ,              0              ,              1              ,              1              ,              0              ,                                                                          0              ,              1              ,              0              ,              0              ,              1              ,              1              ,              1              ,              0              ,              1              ,              1              ,              0              ,              1              ,              1              ,              1              ,              1              ,              0              ,              0              ,              0              ,              1              ,              1              ,              1              ,              1              ,              1              ,              0              ,              1              ,              0              ,              0              ,              0              ,              0              ,              0              ,              0              ,              1                                          }                  L      ⁢                          ⁢      4        =          {                                                  1              ,              0              ,              0              ,              1              ,              0              ,              1              ,              1              ,              1              ,              0              ,              1              ,              1              ,              1              ,              1              ,              0              ,              1              ,              0              ,              0              ,              1              ,              1              ,              0              ,              1              ,              0              ,              1              ,              0              ,              1              ,              1              ,              0              ,              1              ,              1              ,              0              ,              0              ,              0              ,                                                                          0              ,              1              ,              1              ,              1              ,              1              ,              1              ,              0              ,              0              ,              1              ,              0              ,              0              ,              0              ,              1              ,              1              ,              0              ,              0              ,              1              ⁢                                                          ,              1              ,              1              ,              0              ,              0              ,              0              ,              1              ,              0              ,              1              ,              0              ,              0              ,              0              ,              0              ,              0              ,              0              ,              1                                          }      \nConventional MoCA nodes store these binary sequences in a lookup table 214 as shown in FIG. 2. To append a preamble to the beginning of a data packet, each preamble component is accessed from the lookup table 214. Each component is π/4 binary phase-shift keying (BPSK) modulated as a time domain signal at modulator 218. For example, a zero in the binary sequence lookup table 214 gets mapped to a 1+j complex constellation point and a one in the binary sequence lookup table gets mapped to a −1−j complex constellation point. The component is multiplied by a scale factor α at multiplier 208 to ensure that the preamble power is scaled to the appropriate value relative to the data portion of the packet. The scaled component is rounded at rounding block 210 to match the preamble to the bit per sample precision of the DAC 212, which is 10-bits. Once the preamble has been generated, switch 216 will change orientations so it will pass the data from IFFT 206 to DAC 212.\nIn the conventional MoCA network with a 50 MHz bandwidth, generating π/4 BPSK modulated preambles in the time domain by accessing the components stored in the preamble sequence lookup table 214 is sufficient to calibrate a receiver. However, as the bandwidth of a MoCA network is constrained below 50 MHz, generating a preamble using time domain modulation is inadequate to accurately represent the spectrum of the data at a receiver because the time domain nature of the preamble generation does not enable spectral shaping. As a result of the preamble's poor spectral representation of the data portion of the packet, the receiver calibration will suffer and data will be lost.\nAccordingly, a method for generating a spectrally compliant preamble in a MoCA network with a reduced bandwidth is desired."}
{"text": "1. Technical Field\nThe present invention relates to a force detection apparatus and a robot.\n2. Related Art\nFor example, as a force detection apparatus that detects a force applied thereto, a force detection apparatus described in Patent Document 1 (JP-A-10-68665) is known. The force detection apparatus described in Patent Document 1 has an attachment plate, an attachment block, a measurement cell sandwiched between the plate and the block, and two preload screws that apply preloads to the measurement cell through the attachment plate and the attachment block. Further, in the force detection apparatus described in Patent Document 1, the two preload screws are arranged in the width direction of the apparatus side by side for reduction of the height of the apparatus.\nHowever, in the force detection apparatus, the preload screws are arranged in the width direction of the apparatus side by side, and it is difficult to reduce the width of the apparatus."}
{"text": "Decabromodiphenylethane is a time-proven flame retardant for use in many flammable macromolecular materials, e.g. thermoplastics, thermosets, cellulosic materials and back coating applications of very high quality.\nGovernmental regulating agencies tend to be moving away from partially brominated analogs and more towards perbrominated compounds as evidenced by the recent EU RoHS (Restriction on Hazardous Substances) directive (2002/95/EC) relating in part to partially brominated diphenyl oxides. Even with the exemption of decabromodiphenyl oxide from RoHS per 2005/717/EC, the regulations have not been clear enough in terms of the acceptable nonabromodiphenyl oxide content in electrical and electronic products. Some end users therefore find it uncomfortable using the commercial decabromodiphenyl oxide in which significant amounts of nonabromodiphenyl oxide exists as impurity. In order to meet the strictest interpretation of RoHS by the end users, a high purity version of decabromodiphenyl oxide is being marketed by Albemarle Corporation. In view of the confusion concerning the presence of small quantities of lower brominated impurities in the flame retardant products, there is thus a need in the market place for high assay perbrominated flame retardants.\nDecabromodiphenylethane is presently sold as a powder derived from the bromination of 1,2-diphenylethane. Among prior processes for effecting such bromination are the bromination processes described in U.S. Pat. Nos. 6,518,468; 6,958,423; 6,603,049; 6,768,033; and 6,974,887. Decabromodiphenylethane has been commercially produced by the assignee of this application for many years using a standard process. Each batch of product was analyzed by a GC procedure. A review of the GC analyses indicated that the average bromine content of over 4000 batches of decabromodiphenylethane product was 97.57 area percent with a 3-sigma precision of ±1.4 area percent. In some cases, the analysis of the product from a given run provided assays of decabromodiphenylethane in the region of about 99 area percent and above, and in some other cases significantly lower GC assays were obtained. The reasons for this variance cannot be established from the information available.\nGas chromatographic analysis of commercial decabromodiphenylethane products available in the marketplace from other manufacturers have, in some cases, also given assays of a decabromodiphenylethane product as high as about 99.6 area percent. In other cases, GC analyses of commercial decabromodiphenylethane products available in the marketplace have indicated the presence of much lower amounts of decabromodiphenylethane in the product. Information on the method by which such high assay products were produced and the purification procedures used, if any, is not generally available to the public.\nFrom at least the standpoint of providing environmentally-friendly process technology, it would be highly desirable if commercially feasible processes could be found that would produce on a consistent basis a decabromodiphenylethane product that comprises at least about 99.0 GC area percent of decabromodiphenylethane (Br10DPE), with the balance consisting essentially of nonabromodiphenylethane (Br9DPE). Such product is hereinafter often referred to in the specification and claims hereof as “high assay decabromodiphenylethane product”. Moreover, this high assay decabromodiphenylethane product is a “reaction-derived” product which term as used herein including the claims, means that the composition of the product is reaction determined and not the result of use of downstream purification techniques, such as recrystallization or chromatography, or like procedures that can affect the chemical composition of the product. Adding water or an aqueous base such as sodium hydroxide to the reaction mixture to inactivate the catalyst, and washing away of non-chemically bound impurities by use of aqueous washes such as with water or dilute aqueous bases are not excluded by the term “reaction-derived”. In other words, the products are directly produced in the synthesis process without use of any subsequent procedure to remove or that removes nonabromodiphenylethane from decabromodiphenylethane."}
{"text": "1. Field of the Invention\nThis invention relates to a method of analysis and of increasing the speed at which the condition of a database is determined. The invention is more particularly related to the determination of appropriate data needed for a comprehensive analysis of specific tables in the database, and to increase the speed at which the analysis for determining the condition of a database table is performed.\n2. Discussion of the Background\nWith the increasing use of computer databases in each of the home, office, and scientific environments, computer users and Database Administrators (DBA\"\"s) are increasingly confronted with the tasks of determining the condition or health of databases and performing various operations to maintain the database, including reorganization or rebuilding of tables stored in the database. Before rebuilding a database table(s), the DBA must first determine whether the condition of the database table(s) requires rebuilding.\nTypically, commercial database products include tools to perform analysis to determine the condition of a database. For example, Oracle database management systems include an analyze function that performs a predetermined function including an SQL call UPDATE_STATISTICS_FOR_TABLE. The analyze function reads a selected database table and generates a set of statistics that reveal the condition of the table.\nFIG. 1 is a flowchart illustrating the basic steps performed by a conventional Database Management System (DBMS) for determining the condition of a database table. At step 10, a Database Administrator (DBA) provides the name of a table to be analyzed. The table name is provided to a DBMS function that returns data about the named table (step 20). The returned data is analyzed resulting in statistical information (step 30) which is then placed (stored) in a location for retrieval for by DBA (step 40).\nThe DBMS function that returns data (DBMS analysis function) is shown in FIG. 2 and includes the steps of building an SQL call (step 50), invoking the SQL call (step 60), performing SQL overhead and administrative functions (step 70), and returning the SQL results (step 80). The SQL results are provided as the return values (statistical information) in step 30 of FIG. 1.\nOnce the SQL results are analyzed, the displayed (or stored) information is utilized by the DBA to determine appropriate action, if needed, to properly maintain the database (rebuild tables, for example). In the example function, UPDATE_STATISTICS_FOR_TABLE, the Oracle DBMS provides table information including No. of Rows, Chained Rows, Blocks allocated, Blocks allocated but free, Blocks on freelist, Average available space in Freelist, Average free space in a Block, and Min, Max, Average Row Length.\nHowever, the provided table information does not provide enough information to give the DBA the most accurate picture of the DBMS table condition. In addition, current methods for retrieving the DBMS table information are slow, being encumbered by overhead associated with standard DBMS procedures and function calls (file reading techniques, and SQL, for example).\nModern database tables contain enormous amounts of data. If a DBA needs to determine the condition of a large database table, analysis techniques based on current methods can take unacceptably large amounts of time during which a database table may not be available to a user. Furthermore, current methods do not provide the DBA with the exact information needed for a proper analysis. Therefore, the DBA may rebuild tables when it is not necessary or where other action may have been more appropriate. All of the above leads to increased down time and maintenance efforts.\nThe present inventor has realized the need for increasing the amount of data gathered to determine the condition of tables within a database and to increase the speed at which the condition is determined. Accordingly, it is an object of the present invention to provide a comprehensive analysis of at least one database table to determine an information set relating to the health of the tables analyzed.\nIt is another object of this invention to provide method to increase the speed at which database tables are analyzed.\nThese and other objects are accomplished by a method for analyzing database tables including the steps of identifying a database table, determining elements corresponding to the condition of the database table based on data contained in the database tables and providing said elements to one of a display for viewing by a user, a storage location, and a mechanism for using said elements. The elements determined include at least one of number of simple rows, number of migrated rows, number of deleted rows, number of blocks initialized but empty, average free space in a block unused, average free released free space, average row length of migrated rows, average row length of chained rows, number of blocks to be read for normal, chained, and migrated rows, measure of density of blocks used, number of blocks exceeding a percentage used threshold, number of blocks less a percentage free threshold, number of chained rows spanning data files, and min max average of chained rows spanning blocks.\nIn addition, a method to increase the speed of analysis is provided including the steps of identifying a table in said database, reading each of blocks associated with said table using direct asynchronous I/O, and performing said analysis based on information contained in said blocks.\nAdditional speed of analysis is provided by parallel processing and bypassing SQL overhead associated with regular RDBMS architectures. The parallel processing includes the steps of dividing the blocks read into plural sets, and providing each set to separate processes to determine information needed for the analysis."}
{"text": "Conventionally, as a linear guide device for guiding a movable body such as a table or the like along a fixed unit such as a bed or the like, there has been known a bearing for linear sliding in which a sliding base carrying a movable body such as a table or the like is moved along a track rail arranged on a fixed unit such as a bed, a saddle or the like, or a ball spline device in which a cylindrical nut member fitted to a spline shaft is moved along the spline shaft via balls.\nAmong them, as a former bearing for linear sliding, there has been known a constitution comprising a track rail having a ball rolling groove, a sliding base having a load rolling groove opposed to the ball rolling groove as well as a ball return hole in parallel therewith, having a direction change path for guiding balls from the load rolling groove to the ball return hole and moved along the track rail, and a number of balls rolling between the sliding base and the track rail while carrying a load and circulating an endless track constituted by the load rolling groove, the ball return hole and the direction change path of the sliding base.\nFurther, as a latter ball spline device, there has been known a constitution comprising a spline shaft having a ball rolling groove along the axial direction and a nut member having a load rolling groove opposed to the ball rolling groove and fitted to the outer periphery of the spline shaft via balls, where the nut member is moved along the spline shaft similarly in accordance with endless circulation of the balls and which is used by transmitting torque mutually between the nut member and the spline shaft.\nAccording to the conventional linear guide device constituted as described above, the endless track of a sliding member such as the sliding base or the nut member is filled with the balls and accordingly, when the sliding member is moved along the track rail or a guide shaft such as the spline shaft, the balls contiguous to each other are circulated in the endless track while colliding with each other or rubbing each other and there poses a problem where the balls are worn at an early stage and the device life is shortened.\nHence, as means for resolving such a problem, there has been proposed a bearing for linear sliding in which a ball connector where a number of balls are aligned and held is integrated to the endless track (Japanese Unexamined Patent Publication No. JP-A-5-52217). As shown by FIG. 25 and FIG. 26, according to such a ball connector 200, spacers 202 are interposed among respective balls 201 contiguous to each other, the spacers 202 are connected by a pair of strip-like connecting members 203 along a direction of arranging the balls by which the balls 201 are connected in a rosary-like shape and the ball connector is fabricated by molding flexible resin by injection molding in which the balls 201 are arranged in a die as cores.\nAccording to the conventional bearing for linear sliding constituted as described above, as shown by FIG. 27, the above-described ball connector 200 is integrated to a ball return hole 205 and direction change paths 206 of a sliding base 204 and circulated in the endless track and in this case, the spacers 202 are interposed among the balls 201 contiguous to each other and accordingly, mutual friction or collision among the balls is prevented and wear of the balls 201 can be prevented as less as possible.\nMeanwhile, in circulating the balls in the endless track, the ball relieved of a load needs to be scooped up from the ball rolling groove of the track rail to the direction change path of the sliding base and in such a scoop-up operation, as shown by FIG. 28, it is preferable to guide the ball 201 in a direction where a ball rolling groove 208 of the track rail 207 and a load rolling groove 209 of the sliding base 204 are opposed to each other (hereinafter, described as ball contact direction). Because when the ball 201 is guided in such a direction, meandering of the ball 201 in the direction change path 206 is prevented and circulation of the ball 201 in the endless track of the sliding base 204 is carried out smoothly.\nMeanwhile, from a standpoint of avoiding magnifying of the sliding base 204, there is a case where formation of the ball return hole 205 in the ball contact direction is not necessarily proper and there is a case where the ball return hole 205 cannot be formed in the ball contact direction for avoiding interference with a tap hole since the tap hole is formed at the sliding base 204 to fix a movable body such as a table or the like.\nAccordingly, in order to promote the degree of freedom of a position of forming the ball return hole 205 in respect of the sliding base 204 while achieving smooth formation of circulation of the balls 201 in the endless track, as shown by FIG. 29, the direction change path 206 for guiding the balls from the load rolling groove 209 to the ball return hole 205 needs to form to bend from the ball contact direction.\nFurther, in the case of considering the inherent role of the direction change path where the direction of rolling the balls is reversed, it is preferable that the swirl radius of the balls in such a direction change path is large and for that purpose, the length of the direction change path needs to set long. However, in the case where the ball return hole is positioned in the ball contact direction in view from the load rolling groove, when the length of the direction change path is set long, an interval between the load rolling groove and the ball return hole is naturally widened and magnification of the sliding base cannot be avoided.\nAccordingly, also in view thereof, the ball return hole needs to form at an arbitrary position of the sliding base and the direction change path for communicating and connecting the ball return hole with the load rolling groove needs to form to bend.\nHowever, according to such a conventional ball connector 200, side edges of the strip-like connecting members 203 for connecting the respective spacers 202 are linearly formed and therefore, there poses a problem in which although the ball connector 200 is easy to bend in a direction intersecting with the faces of the connecting members 203 (arrow line A direction in FIG. 26), the ball connector 200 is difficult to bend in a direction in parallel with the faces of the connecting members 203 (arrow line B direction in FIG. 25), when the direction change path 206 is bent to deviate from the ball contact direction as mentioned above, the ball connector 200 is obliged to integrate into the endless track by forcibly bending it and therefore, smooth circulation of the ball connector 200 cannot be expected.\nTherefore, according to the conventional linear guide device using the ball connector, when smoothness of circulation of the ball connector in the endless track is considered, the ball return hole is obliged to form in the ball contact direction in view from the load rolling groove which constitutes one factor for magnifying the sliding base.\nFurther, according to the conventional ball spline device described above, from a standpoint of achieving to downsize the nut member, the ball return hole needs to install as near to the inner diameter of the nut member as possible and it is difficult to form the ball return hole in the ball contact direction in view from the load rolling groove similar to the bearing for linear sliding mentioned above.\nTherefore, when the conventional ball chain shown by FIG. 25 is integrated to the endless track of the ball spline device, smooth circulation of the ball chain cannot be expected.\nMeanwhile, according to a ball screw device which is used with a purpose of converting rotational movement of a motor into linear movement and which constitutes a linear guide unit of a machine tool or the like by being used along with the linear guide device mentioned above, a screw shaft and a nut member are in mesh with each other via balls and the nut member is similarly provided with an endless track for the balls. Accordingly, from a standpoint of preventing wear of the balls circulating at inside of the nut member, it is preferable also for the ball screw device to integrate a ball connector to the endless track. However, as already has been explained, the conventional ball connector is provided with a structure where it is easy to bend only in specific directions and accordingly, it is difficult to integrate the ball connector to the ball endless track of the ball screw device.\nThat is, in integrating the ball connector to the ball endless track of the ball screw device, the ball connector needs to wind spirally in respect of the screw shaft. Accordingly, the ball connector needs to circulate not planarly but sterically and in the case of the conventional ball connector having small degree of freedom in respect of the bending direction, smooth circulation of the ball connector in the endless track cannot be expected. Further, when the ball connector is integrated to the ball screw device, the ball connector is obliged to circulate in the endless track in a state accompanied by more or less twist and when the conventional ball connector having a low degree of freedom in respect of the bending performance is used by forcibly twisting it, there is a concern where the balls are detached from the ball connector per se."}
{"text": "1. Field of the Invention\nThe present invention relates to electrical cable connectors, and more particularly to a connector for electroluminescent cable having coaxial conductors.\n2. Description of the Related Art\nElectroluminescent cable (EL-cable) is a cool to the touch, bendable, vinyl coated wire that emits a pleasant 360-degree softly glowing neon light. EL-cable is a flexible wire cable having a solid copper center conductor surrounded by a material which is luminescent in an electric field. Two thin filaments or wires which are shorted together and helically wound around the luminescent material. The assembly is covered with one or two layers of vinyl or other plastic insulating material.\nWhen an alternating current is conducted through the center conductor and the two filaments, the alternating electromagnetic field between the conductors causes the luminescent material to glow. Although the EL cable may be powered directly from the A.C. power mains, frequently the cable is powered by a D.C. inverter connected to a battery. The color emitted by the cable may vary with the frequency of the A.C. voltage or current. Usually the voltage must exceed a minimum threshold voltage before the EL-cable will glow.\nEL technology is relatively new and only within the past few years has EL-cable become available in consumer products, specifically applications requiring lengths of glowing lights, applications which previously employed, LED or other lamp technologies. An efficient and effective method for connecting EL-cable either to other strands of EL-cable or to a pair of copper wires has not been adequately addressed.\nIn general, connectors for electric cables are not new and the technology is well represented by devices for splicing wires together and for connecting wires to electronic devices. U.S. Pat. No. 4,921,451, issued to R. Carlson in May of 1990, discloses in-line fuse holders for two-bladed fuses which can be fastened in series to an electrical wire by severing the wire in which the holder is to be incorporated, inserting the severed ends of the wire into the holder, and mechanically fastening the wire securely in the holder.\nU.S. Pat. No. 5,007,855, issued to O'Brien et al. in 1991, discloses a cable connector having a pair of electrically conductive jumper elements with a pair of spaced sharp protrusions that are electrically connected. U.S. Pat. No. 5,055,071, issued to Carlson, deceased et al. in October of 1991, describes a cable connector in which two cables' conductors are each engaged by a slotted conductor, both of which engage with a common conducting bridge.\nU.S. Pat. No. 5,702,262, issued to Brown et al. in December of 1997, discloses a housing having connectors in coaxial alignment with a pair of barrels. In U.S. patent application Publication No. 2002/0182934, published in December 2002, Endo et al. describes a coaxial connector having a central contact, an insulating housing, a grounding shell and a clamp. A crimp barrel serves as a conductor-connecting portion that is crimped into contact with the central conductor of the coaxial cable once the central conductor is inserted into the central contact.\nNone of the above inventions and patents, taken either singly or in combination, is seen to describe a connector for EL-cable as claimed. Thus a simplified electroluminescent cable connector solving the aforementioned problems is desired."}
{"text": "The present invention relates to a method for producing a standard color model used for judging the shade and tone of beef-color.\nUp to the present, the market standard of beef carcass is chiefly based on its quantity, because a simple and exact method for judging its quality has not yet been established. Accordingly, the grading of market beef carcass depends on a subjective decision of a judge. Therefore, it is necessary to provide a simple, accurate and exact method for judging the quality of meat by colorimetrically comparing the tone thereof.\nThe tone of beef is an important factor for determining its quality. Means for judging the quality of chemical analysis or by using a color comparator or only with human eyes are already known. However, these means involve either complex operations or expensive measuring tools. Additionally judgements were often arbitrary."}
{"text": "Many techniques have been employed for the treatment of impure water so as to provide a potable water supply. Generally the most widely employed techniques involving minimum expense utilize distillation of one form or another, with wide variations in the source of heat energy being employed. One of the major impediments to use of distillation involving rapid boiling of the water is the tendency for foaming to occur with certain types of contamination. This foaming often will require excessively large dimensions for the boiler or an intricate and elongated passage from the boiler to the condenser. Moreover, even with water feedstocks which do not evidence substantial foaming, rapid boiling will frequently produce a tendency for water droplets which are carried by the steam vapors into the condenser. As will be readily appreciated, foam droplets or such water droplets will exhibit the contamination of the feedstock and thereby contaminate the condensate which is produced.\nAnother problem which is frequently encountered in certain sections of the world involves the contamination of the water supply by organic compounds which vaporize along with the water in the boiler. Exemplary of such contaminants are carbon tetrachloride and chloroform which result from chlorination of water supplies containing dissolved methane and other hydrocarbons.\nVarious types of water purification devices have been suggested for providing potable water in the home, or on shipboard, or in other locations. Many of these devices require substantially continuous operation for efficiency, others require substantial investments, and still others require extensive maintenance from time to time.\nIt is an object of the present invention to provide a novel water purification apparatus which substantially eliminates carry-over contamination of the condensate.\nIt is also an object to provide such an apparatus which may be simply and ruggedly constructed to provide a durable and readily serviceable unit.\nAnother object is to provide such apparatus which is capable of substantially eliminating dissolved gaseous contaminants.\nStill another object is to provide such apparatus which may be discontinuously operated, and which, during operation automatically, regulates flow of water into the boiler and which, upon discontinuance of operation, automatically effects draining of the several components.\nA further object is to provide a method whereby water feedstock containing impurities may be rapidly and efficiently distilled without carry-over contamination."}
{"text": "The invention relates generally to the field of methods and systems for interactive computing. The invention relates more particularly to the fields of methods for optimization.\nGenetic and evolutionary algorithms (GAs) may be employed to optimize a solution set in a variety of problem types. Generally, GAs include selection from among a plurality of possible solutions based on fitness and recombination of one or more selected solutions to generate additional possible solutions. These two actions may be repeated as necessary until there is a convergence to an optimal or suitable solution.\nAn important concept used in GAs to implement a survival-of-the-fittest mechanism and to distinguish between good and bad solutions is the notion of fitness. Unlike traditional search methods, a fitness measure can be relative, and it can be an objective measure including a mathematical equation, model, or a computation. It could also be a subjective measure involving human evaluation, or even be coevolved in a co-operative or competitive environment.\nGAs may be classified generally based on whether selection (choosing) and recombination (innovation), respectively, are performed by a human or by a computational agent such as, but not limited to, a computer. For example, a standard GA relies on computational selection (e.g., via computing a fitness function, model, or computation) and recombination (e.g., crossover and/or mutation). Computer-aided design (CAD), on the other hand, relies on computational selection of solutions, but a human performs the recombination of the solutions. In a human-based GA, a human performs both steps, though these steps may be augmented, for example, by a network or other communication or collaborative medium.\nGAs in which fitness measure or quality of candidate solutions (e.g., selection or choosing) is provided by human evaluation and/or judgment rather than a fitness function, but in which recombination (e.g., innovation) is at least in part based on a computational agent are referred to as interactive genetic algorithms (iGAs). iGAs, a method of interactive evolutionary computation, allow human interaction to provide subjective input regarding fitness, while still providing computational recombination.\nMore particularly, in conventional iGAs, qualitative fitness as determined by a human user replaces quantitative fitness. As one example of providing qualitative fitness, a subset of generated solutions may be sorted and presented to a user, who selects a preferred solution from among the presented ones. The preferred solution is used for computationally-provided recombination of new subsets, which are again presented to the user.\nDawkin's Blind Watchmaker program, described in Dawkins, R., The Blind Watchmaker, New York: W. W. Norton, 1986, and the Faceprints system developed at New Mexico State University, as described in Caldwell, C., and Johnston, V. S., “Tracking a criminal suspect through face-space with a genetic algorithm”, Proceedings of the Fourth International Conference on Genetic Algorithms, Morton Kaufmann, 1991, pp, 416-421, are two nonlimiting examples of iGAs. For example, in Faceprints, a system replaces the role of a human sketch artist in evolving the faces of criminal suspects from witness recollection. Faces are encoded as binary strings, where subcodes represent different facial features (nose type, mouth type, hair type, etc.). Each full chromosome maps to a face, and the population of chromosomes is presented to the human critic, who is asked to determine how closely the face resembles that of the criminal. This subjective ten-point scale is used to drive the evolution of subsequent generations of faces, and in a relatively short time, the iGA arrives at a reasonable facsimile of the correct face.\nHowever, unlike in evolutionary algorithms with objective fitness measures, one of the daunting challenges of iGAs is providing effective methods of combating user fatigue. Because a human is in charge of the evaluation of a solution, user fatigue results from the large time scale between user evaluations and the evolutionary mechanisms used to generate new solutions until convergence. For example, even for moderately sized problems, iGAs may require, for example, a few hundred to a few thousand fitness evaluations, which is highly improbable—sometimes even impossible—for users to perform. User fatigue may occur at fairly short time periods, such as within 1-2 hours, or even earlier.\nUser fatigue can result in sub-optimal solutions, or the end of an iGA session before even a suitable solution is found. Similarly, repeated evaluation of similar solutions can result in user frustration relatively quickly, compounding user fatigue.\nAdditionally, an important element for the successful application of interactive genetic algorithms (iGAs) or any other interactive evolutionary computation method is the reliability of the evaluations, or decisions, provided by a user. The decisions that the user makes when using an iGA guide the search across the space of possible hypotheses. Any interactive procedure, regardless of how efficient it can be for reducing the number of evaluations required, will fail to provide high-quality solutions if the user is unable to provide consistent fitness evaluations.\nThus, it is desired to provide ways to improve efficiency enhancement for iGAs, while also providing the reliability needed to use iGAs successfully."}
{"text": "1. Field of the Invention\nThe present invention relates to non-volatile memory cells. More particularly, the present invention relates to push-pull non-volatile memory cells and to arrays of such memory cells.\n2. The Prior Art\nPush-pull non-volatile memory cells are known in the art. FIG. 1 is a schematic diagram shows an illustrative prior-art push-pull non-volatile memory cell. The memory cell includes a p-channel non-volatile memory transistor connected in series with an n-channel non-volatile transistor. The p-channel non-volatile memory transistor and the n-channel non-volatile transistor may be fabricated as floating-gate flash transistors or may be fabricated using other known non-volatile transistor technologies.\nThe memory cell shown in FIG. 1 may be used to drive a switch transistor such as the n-channel transistor shown with its gate coupled to the common drain connections of the p-channel non-volatile memory transistor and the n-channel non-volatile transistor. Such a memory cell arrangement may be used to form programmable circuit connections in a programmable integrated circuit such as a field programmable gate array (FPGA) integrated circuit.\nThe memory cell shown in FIG. 1 may be programmed into one of two states. In the first state, the n-channel non-volatile transistor is turned on and the p-channel non-volatile transistor is turned off. In this state, the gate of the n-channel switch transistor is grounded, turning it off. In the second state, the p-channel non-volatile transistor is turned on and the n-channel non-volatile transistor is turned off. In this state, the gate of the n-channel switch transistor is at approximately VDD, turning it on.\nIn the case of flash non-volatile memory transistors, the p-channel non-volatile transistor is programmed by placing a negative voltage, such as −4v, on the source of the p-channel non-volatile transistor, placing a positive voltage such as 8.5v on its gate, while its bulk is biased at a voltage such as 1.2v. The gate and source of the p-channel non-volatile transistor are biased at 0v during this procedure.\nThe n-channel flash non-volatile transistor is programmed by placing a voltage such as 8.5v, on its gate, and placing a voltage such as 4.5v on its source. The gate of the p-channel non-volatile transistor is biased at a voltage such as −3.3v and its source is biased at a voltage such as 0v during this procedure."}
{"text": "1. Field of the Invention\nThe present invention relates to composite materials technology, and more specifically to a relatively light-weight, inexpensive, durable, high performance structural laminate composite material for use to 1000xc2x0 F., and above, which can advantageously be used in high temperature environments. More particularly, the preferred embodiment of the present invention relates to a graphite-fiber/phenolic-resin composite material which retains relatively high strength and modulus of elasticity at temperatures as high as 1,000xc2x0 F. (538xc2x0 C.). The material costs only 5 to 20 percent as much as refractory materials do. The fabrication of the composite includes a curing process in which the application of full autoclave pressure is delayed until after the phenolic resin gels. This modified curing process allows moisture to escape, so that when the composite is subsequently heated in service, there will be much less expansion of absorbed moisture and thus much less of a tendency toward delamination. In contrast, internal pressure caused by the expansion of moisture absorbed in other prior art composite materials like prior art graphite/epoxies and prior art graphite/polyimides causes delamination at temperatures in the range of 500 to 700xc2x0 F. (260 to 370xc2x0 C.).\n2. General Background\nAt the request of NASA/MSFC, Martin Marietta Manned Space Systems has performed an extensive development/verification activity for a composite nose cone for the external tank (ET). At the time of the initiation of this effort, there was no materials technology available to provide a nose cone which could withstand the high heating and structural loading of the ET nose cone without (a) requiring the use of secondary heat shield materials, (b) increasing the weight of the existing nose cone, and (c) significantly increasing the cost over the existing nose cone cost. There were high temperature polymeric composite materials available; however, none met all requirements. Carbon/phenolic laminates have been proven in rocket nozzle applications to be able to withstand extreme heating conditions; however, these materials did not possess the specific strength and stiffness required for a weight-effective structure. Also, recent data shows that the materials on the market today have the potential to xe2x80x9cply lift,xe2x80x9d or delaminate due to internal pressure caused by absorbed moisture, at about 500xc2x0 F. Graphite/polyimide laminates showed promising mechanical properties, but suffered from the moisture-induced delamination problem (also known as xe2x80x9cthermal shockxe2x80x9d) at temperatures below 700xc2x0 F. in laminates of the thickness required for a composite nose cone. Other technologies such as ceramic matrix composites and carbon/carbon were considered too expensive for this application. Therefore, a program was initiated to develop laminate material which could meet all requirements.\nU.S. Pat. No. 3,724,386 for xe2x80x9cAblative Nose Tips and Method for their Manufacturexe2x80x9d discloses in Example II heating graphite yarn impregnated with phenolic resin slowly to 160xc2x0 F. to slowly evaporate solvent from the resin (see column 8, lines 16-18).\nU.S. Pat. Nos. 4,100,322 and 4,215,161 for xe2x80x9cFiber-Resin-Carbon Composites and Method of Fabricationxe2x80x9d disclose impregnating graphite yarn with phenolic resin under vacuum and a temperature of about 150xc2x0 F. until the solvent has gone and the resin gels, then further heating the composite to cure it. However, the solvent stripping process was interrupted twice and each time pressure of 200 psig was applied to the composite material. It is then subjected to pyrolysis, and then pores of the composite are impregnated with phenolic resin. After this, the phenolic resin is cured at about 350xc2x0 F. The resulting structure is said to be graphite/carbon/phenolic composite, and its porosity is disclosed to be 4%. A carbon/carbon/phenolic composite described therein is said to have a porosity of 5.8%.\nU.S. Pat. No. 4,659,624 for xe2x80x9cHybrid and Unidirectional Carbon-Carbon Fiber Reinforced Laminate Compositesxe2x80x9d discloses a method similar to the method disclosed in U.S. Pat. Nos. 4,100,322 and 4,215,161 (and with similar materials), but one in which more resin is added and pyrolized up to 5 times. This patent points out at column 2, line 50 through column 3, line 2 that it is important to properly initially cure laminate materials to provide interconnecting pores which allow the escape of gases formed during post-cure pyrolysis.\nU.S. Pat. No. 4,957,801 for xe2x80x9cAdvance Composites with Thermoplastic Particles at the Interface Between Layersxe2x80x9d discloses a resin-impregnated fiber layer with outer layers of resin thereon. The fiber can comprise, for example, graphite.\nU.S. Pat. No. 5,288,547 for xe2x80x9cToughened Resins and Compositesxe2x80x9d discloses a composite in which a porous membrane film of thermoplastic material is sandwiched between two layers of resin-impregnated fibers, and then the composite is cured in an autoclave, for example. The resin can be, for example, phenolic resin.\nU.S. Pat. No. 5,359,850 for xe2x80x9cSelf Venting Carbon or Graphite Phenolic Ablativesxe2x80x9d discloses a resin-impregnated reinforcing cloth made of, for example, graphite fibers with degradable fibers interwoven therewith. The degradable fibers are chosen such that they degrade at a temperature of about 400xc2x0 F. to 500xc2x0 F. so that they will provide passageways for the gaseous decomposition products produced as the resin matrix approaches the char temperature. In this patent, foreign material is introduced to create porosity. The fabric weave is altered by introducing a low-temperature degradable thread which may not assure fabric strength properties. The porosity which is created by this process is uniform. There is a definite pattern when the foreign material is replaced by voids. It is believed that the addition of these special degradable fibers will add to the cost of the material. Further, it is believed that in some cases the degradable fibers might not burn away before the plies blow apart.\nU.S. Pat. No. 5,360,500 for xe2x80x9cMethod of Producing Light-Weight High-Strength Stiff Panesxe2x80x9d discloses a panel made by a pair of surface members separated and supported by an internal core in which spaces or interconnected pores provide vents to an edge of the panel so that gas can flow through the vents during a pyrolysis process. The vents are on the order of 10 mm in diameter.\nNone of these patents discloses a composite material with a weight, thickness, structural performance, and pore structure as advantageous for use in a nose cone of the external tank of the space shuttle, or other high temperature structural applications, as the material of the present invention.\nA novel materials technology has been developed and demonstrated for providing a high modulus composite material for use to 1000xc2x0 F. The material of the present invention can be produced at 5-20% of the cost of refractory materials, and has higher structural properties. This technology successfully resolves the problem of xe2x80x9cthermal shockxe2x80x9d or xe2x80x9cply lift,xe2x80x9d which limits traditional high temperature laminates (such as graphite/polyimide and graphite/phenolic) to temperatures of 550-650xc2x0 F. in thicker (0.25xe2x80x3 and above) laminates. The technology disclosed herein is an enabling technology for the nose for the External Tank (ET) of the Space Shuttle, and has been shown to be capable of withstanding the severe environments encountered by the nose cone through wind tunnel testing, high temperature subcomponent testing, and full scale structural, dynamic, acoustic, and damage tolerance testing.\nIn the present invention, cure conditions (temperature, pressure, vacuum) and cure apparatus (specific vacuum bag methodology) are manipulated to produce a graphite/phenolic composite laminate with a permeable microstructure comprising an interconnected network of pores which allows moisture to escape from the composite material when the composite material is heated; this helps prevent delamination (xe2x80x9cply liftxe2x80x9d or xe2x80x9cthermal shockxe2x80x9d) when the material is heated to temperatures above 500xc2x0 F. The graphite/phenolic composite of the present invention can be used for components for applications requiring high strength and stiffness upon exposure to very high heating (e.g. rocket nozzles for missiles or launch boosters, fire walls, heat shields, circuit boards, secondary structure on missiles or launch vehicles which see high aerodynamic heating, and parts to be used on the leading edge of aerodynamic products (airplanes, jets, rockets, fuel tanks for aerospace structures, etc.)).\nThe present invention comprises a method of producing a composite material, comprising the steps of:\nimpregnating a fiber material with a resin to create a resin-impregnated fiber material;\nwithout applying pressure, heating the resin-impregnated fiber material under vacuum at a sufficient temperature for a sufficient amount of time until the resin gels; and\napplying temperature (and, optionally, pressure) for a sufficient period to cure the resin-impregnated fiber material. The starting percentage by weight of fiber material (before being cured) is preferably 30-80%, with the balance resin. The resulting porosity of the composite material is preferably at least 3% by volume, more preferably about 3-25% by volume, and most preferably about 7-14% by volume.\nThe preferred embodiment of the method of the present invention of producing a composite material comprises the steps of:\n(i) impregnating a graphite fiber material with a phenolic resin to create a resin-impregnated fiber material, in a ratio of 30-80% by weight graphite fiber and 20-70% by weight phenolic resin;\n(ii) placing the resin-impregnated fiber in an autoclave or oven;\n(iii) applying full vacuum and/or pressure;\n(iv) raising the temperature to cause the resin to flow and initiate cure,\n(v) holding the material at a temperature to allow gellation of resin while volatiles are being released;\n(vi) raising the temperature for final cure if required;\n(vii) cooling the material;\n(viii) removing the material from the autoclave or oven;\n(ix) post-curing the composite laminate material removed from the autoclave, if required.\nThe present invention includes the composite material made by the method of the present invention disclosed herein, as well as a composite material, produced by any method, having a composition and structure which is the same as the composite material produced by the method of the present invention disclosed herein.\nThe material of the present invention comprises a high performance structural laminate composite material for use in high temperature applications, consisting essentially of resin-impregnated fiber, the resin-impregnated fiber consisting essentially of:\n(a) preferably 50-80% by weight fiber, and\n(b) preferably 20-50% by weight cured resin, the composite material having:\n(c) a permeability sufficient to allow moisture to escape from the composite material, without causing plylift, when the composite material is heated to temperatures up to 1000xc2x0 F. More preferably, the permeability is sufficient to allow moisture to escape from the composite material, without causing plylift, even when the composite material is heated to temperatures above 1000xc2x0 F. The material of the present invention has a microscopic construction which provides permeability that is sufficient to allow moisture to escape therefrom as it is heated to temperatures up to 1000xc2x0 F. and above without exhibiting ply-lift.\nThe composite material preferably has an across-ply permeability having a Darcys constant of at least 10xe2x88x9215 cm2. More preferably, the across-ply permeability of the composite material has a Darcy\"\"s constant of at least 10xe2x88x9214 cm2. Most preferably, the across-ply permeability of the composite material has a Darcy\"\"s constant of at least 10xe2x88x9213 cm2.\nThe material of the present invention comprises a high performance structural laminate composite material for use in high temperature applications, consisting essentially of phenolic resin-impregnated graphite fiber, the phenolic resin-impregnated graphite fiber consisting essentially of:\n(a) preferably 50-80% by weight graphite fiber; and\n(b) preferably 20-50% by weight cured phenolic resin, the composite material having:\n(c) a permeability sufficient to provide a network of pores which allows moisture to escape from the composite material, without causing plylift, when the composite material is heated.\nThe percentage by weight of graphite fiber is more preferably 60-80%, and the percentage by weight of cured phenolic resin is more preferably 20-40%. Most preferably, the percentage by weight of graphite fiber cloth is 65-75%, and the percentage by weight of cured phenolic resin is 25-35%.\nPreferably, the porosity is 3-25% by volume. Most preferably, the porosity is 7-14% by volume.\nPreferably, the compressive strength of the material after exposure to temperatures above 700xc2x0 F. for several minutes is at least 50% of the compressive strength of the material immediately after being cured, the shear strength of the material after exposure to temperatures above 700xc2x0 F. for several minutes is at least 50% of the shear strength of the material immediately after being cured, and the compressive strength of the material at 900xc2x0 F. is at least 25% of the compressive strength of the material at room temperature.\nThe graphite fiber cloth was selected to have a combination of high strength, high modulus, good thermo-oxidative stability, and moderate cost. The optimum fiber type to provide this balance is a fiber made from a polyacrylynitrile (PAN) precursor, such as the Toho G30-5001 fiber used in the development documented herein. Similar fibers are Hercules AS4 and IM-7, and Amoco T300, T650-35, and T650-45. Fiber types which were not selected were fibers based on pitch precursors (e.g., Amoco P-75 and P-100), or fibers based on rayon precursors. Pitch based fibers are much more expensive and do not have adequate strength. Rayon based fibers do not have the desired strength or modulus. The selected fiber was woven into an 8-harness satin fabric to facilitate part fabrication. The selected fiber can be, for example, an eight-harness fabric woven from Toho G-30/500-3K graphite fiber. The resin is advantageously selected from a group consisting of phenolics, bismaleimides (BMIs), polyimides, cyanate esters, epoxies, or any blend of these resins. The resin can comprise Cytec 506 phenolic resin.\nPreferably, the graphite fiber material has a minimum tensile strength of at least 300 KSI, more preferably at least 400 KSI, and most preferably at least 500 KSI, a minimum modulus of at least 20 MSI, more preferably at least 25 MSI, and most preferably at least 30 MSI, and relatively low cost. Most preferably, the resin is phenolic resin.\nThe material can consist of phenolic resin-impregnated graphite fiber cloth, and the phenolic resin-impregnated graphite fiber cloth can consist of graphite fiber cloth and cured phenolic resin.\nThe present invention also includes apparatus comprising a component which requires high strength and stiffness upon short term exposure to very high heating, made of the material of the present invention. The component can be a rocket nozzle, a part for an aerodynamic vehicle, or some other component exposed to high heating. The component can be part of a fire wall or heat shield.\nFurther, the present invention comprises vacuum bag apparatus for producing a composite laminate material having a network of pores. This vacuum bag apparatus can comprise:\n(a) a base for receiving the laminate material thereon;\n(b) a non-stick layer to be received on the laminate material for helping to prevent the laminate material from sticking to layers above the non-stick layer;\n(c) a first volatiles flow and resin retaining layer above the non-stick layer for allowing volatiles, but not the majority of the resin, to escape from the laminate material through the first volatiles flow and resin retaining layer as heat is applied and the vacuum is drawn in the bag apparatus;\n(d) a bleeder layer on the first volatiles flow and resin retaining layer for absorbing most of the resin which flows through the first volatiles flow and resin retaining layer;\n(e) a second volatiles flow and resin retaining layer on the bleeder layer for allowing volatiles, but very little resin, to flow through the bleeder layer as heat is applied and the vacuum is drawn in the vacuum bag apparatus;\n(f) a first gas-flow layer on the second volatiles flow and resin retaining layer for allowing gas to flow evenly through the vacuum bag apparatus when a vacuum is drawn in the apparatus;\n(g) a lateral gas-flow layer surrounding the laminate material to ensure that volatiles can flow out of the laminate in virtually any direction;\n(h) a vacuum bag layer attached to the base in an air-tight manner, the base and the vacuum bag layer enclosing the laminate material and the non-stick layer, the first volatiles flow and resin retaining layer, the bleeder layer, the second volatiles flow and resin retaining layer, the gas-flow layer, and the lateral gas-flow layer. In certain circumstances, one or more of the layers can be omitted, as described further below. The vacuum bag apparatus preferably also comprises a port in the vacuum bag layer communicating with a vacuum source for allowing a vacuum to be pulled in the bag. In a room at standard temperature and pressure, the vacuum causes a pressure of about 15 psi to be applied to the laminate in the vacuum bag. The application of additional pressure may not be a necessary step to make the present invention work.\nIt is an object of the present invention to provide a high-strength, low weight, high temperature material which has sufficient permeability to allow moisture to exit therefrom, even when the material has a thickness of more than 0.40xe2x80x3, when heated to temperatures of above 500xc2x0 F., without damaging the internal structure of the material.\nIt is an object of the present invention to provide a high-strength, low weight, high temperature material which has sufficient permeability to allow moisture to exit therefrom, even when the material has a thickness of more than 0.40xe2x80x3, when heated to temperatures of above 1000xc2x0 F., without damaging the internal structure of the material.\nIt is another object of the present invention to provide a method of making such material.\nA further object of the present invention is to provide components made of such material.\nIt is also an object of the present invention to provide a material which can withstand the high heating and structural loading of the ET nose cone without (a) requiring the use of secondary heat shield materials, (b) increasing the weight of the existing nose cone, and (c) significantly increasing the cost over the existing nose cone cost.\nAnother object of the present invention is to provide an ET nose cone made of this material.\nUnlike many prior art methods of producing composite material, in the method of the present invention, there is no pyrolizing step (the composite material of the present invention is not pyrolized). In the method of the present invention, unlike the method of U.S. Pat. No. 5,359,850: no foreign material is introduced to create porosity; the fabric weave is not altered and areal weight of fabric is constant, which assures strength properties; the porosity which is created by the process of the present invention is random and spread over the composite laminate; no material is decomposed by the method of the present invention; and the cost of creating the porosity is relatively low.\nBecause of the high permeability of the material of the present invention, it is believed by the inventors that there will be no ply lift at any thickness, whether the laminate is at least 0.1 inch thick, at least 0.2 inch thick, at least 0.4 inch thick, or even more than 4 inches thick.\nAlthough the specific examples described herein relate to graphite fiber and phenolic resin, other appropriate fibers and resins could be used in conjunction with the present invention."}
{"text": "Commonly owned U.S. Pat. Nos. 6,015,062; 7,069,763; 7,370,507 and published application US2006-0011633 and disclose containers/cans and caps therefor, tooling for making the can bodies and caps, and the method or processes of making cans that employ a reclosable cap. The disclosure of each of these U.S. patents and published applications is incorporated herein by reference. Heretofore, this technology has been used to manufacture new can bodies and oftentimes includes the manufacture of a separate dome that is connected or seamed to the open end of a main can body. The dome in an exemplary arrangement typically has a conical tapering conformation that narrows to an opening that receives a removable cap. To retain the cap in closed relation with the opening, external thread portions are provided adjacent one end of the dome about an outer periphery of the opening. These thread portions are selectively engaged by thread lugs that extend from an inner peripheral portion of a cap rim to draw a cap into sealed, reclosable engagement with the periphery of the opening. In other instances, the thread lugs are formed on an outer peripheral portion adjacent the open end of a can body which may not be tapered. The thread portions are still adapted to receive the corresponding thread lugs of a reclosable cap.\nIn the food and beverage can industry, for example, open ended can bodies having a bottom wall that is integrally formed with the side wall. A blank of material is punched from a cold-rolled sheet and initially formed into a cup that has an end with integral sidewall. The cup is subsequently formed or ironed with dies and forming machinery to shape the metal and form a can body into an open-top can. An open or upper, second end is subsequently closed with a separately formed lid or end panel. The end is connected along a peripheral portion to the open-ended sidewall of the can body via a crimp or seam. Typically, the ends are seamed to the open end of the can body and the consumer accesses the contents of the can by (i) removing the end with a can opener, (ii) removing or tearing away the end via a pull tab, or (iii) an easy-open, retained scored region or panel portion in the end that is opened with a retained tab.\nThere are times, however, where the entire contents of the food or beverage container, for example, may not be used. The consumer must then empty the entire contents from the can since there is no effective way to close and seal the can for future use once the can has been opened. Likewise, these cans are made time and again at the same volume and do not easily lend themselves to converting to new sizes of containers without a significant capital investment in equipment.\nTherefore, a need exists for providing a reclosable can, particularly in the metal food container industry. There is also a desire to accomplish this goal without adding additional metal to the container, while advantageously providing additional volume."}
{"text": "The neuronal cell membrane comprises a phospholipid bilayer and numerous embedded ion channels traversing the membrane. Embedded proteins transport charged molecules altering the electrochemical gradient across the membrane. In response to a stimulus, a voltage differential across the membrane crosses a depolarization threshold, resulting in an action potential.\nNumerous studies have established that the neuronal cell membrane displays the electrical characteristics of resistance (R), derived from embedded ion channels, and capacitance (C), derived from the phospholipid bilayer structure, as a parallel-RC circuit. The parallel-RC circuit equivalent circuit of all cell membranes, including neurons, uses the resistance and capacitance of the phospholipid bilayer. Nerves are preferential electrical current conduction pathways through tissue, and their presence in living tissue makes the tissue anisotropic to electrical current flow. This anisotropicity and phenomenological inductance of the neuronal cell membrane may be demonstrated using externally applied fields.\nThe impedance neurography developments based on this anisotropicity of nerve conduction pathways, in combination with methods to develop a transmembrane current flow, have led to insights regarding the neuronal cell membrane electrical responses. As a result, a wide range of clinical applications may be advanced."}
{"text": "In many commercial endeavors, a plurality of computer applications share a set of data. For example, in the field of medicine, a user can input the same information describing a particular patient to multiple computer applications. The information can be patient information (such as name, sex, social security number, primary physician, medical history, drug interactions, etc.), user information (such as user identification, password, etc.), encounter or visit information (such as date, encounter number, name of treating physician, etc.), observation information (such as diagnosis, blood test results, etc.), financial information (such as insurance coverage, billing history, etc.), or other types of information. The user's task historically entailed repeatedly entering the same data describing the patient into multiple computer applications. This practice of repeatedly entering the same data for multiple computer applications not only applies to the medical field, but also extends to other fields.\nThe healthcare industry recognizes the need for managing and sharing data, and in 1999, the Health Level Seven (HL7) published the first Clinical Context Object Workgroup (CCOW) standard (hereafter “the CCOW standard”) for context management. The CCOW standard defines a context management architecture (CMA) and processes for managing information describing a subject across a range of clinical and other healthcare-related computer applications. The nature of the CCOW standard is described, for example, in HL7 Context Management Standard, Version 1.5, which was rectified in May 2004.\nWhile the CCOW standard allows context and data sharing, a user still must repeatedly enter the same information into multiple computer applications. For example, following a patient visit, a nurse may have to use the same login information and enter the same patient identification information into three or four different computer applications. Repeatedly entering the same information into multiple applications can be burdensome and time consuming.\nAccordingly, a need exists for a method and a system that allows users to easily repeat tasks to streamline workflow and that also allows users to reduce the time wasted on redundant tasks.\nFor simplicity and clarity of illustration, the drawing figures illustrate the general manner of construction, and descriptions and details of well-known features and techniques may be omitted to avoid unnecessarily obscuring the invention. Additionally, elements in the drawing figures are not necessarily drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help improve understanding of embodiments of the present invention. The same reference numerals in different figures denote the same elements.\nThe terms “first,” “second,” “third,” “fourth,” and the like in the description and in the claims, if any, are used for distinguishing between similar elements and not necessarily for describing a particular sequential or chronological order. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the invention described herein are, for example, capable of operation in sequences other than those illustrated or otherwise described herein. Furthermore, the terms “comprise,” “include,” “have,” and any variations thereof, are intended to cover a non-exclusive inclusion, such that a process, method, article, or apparatus that comprises a list of elements is not necessarily limited to those elements, but may include other elements not expressly listed or inherent to such process, method, article, or apparatus.\nThe terms “left,” “right,” “front,” “back,” “top,” “bottom,” “over,” “under,” and the like in the description and in the claims, if any, are used for descriptive purposes and not necessarily for describing permanent relative positions. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the invention described herein are, for example, capable of operation in other orientations than those illustrated or otherwise described herein. The term “coupled,” as used herein, is defined as directly or indirectly coupled in an electrical or non-electrical manner."}
{"text": "Field of the Invention and Related Art Statement\nThe present invention relates to a method of effecting an immunological analysis.\nDue to recent progress in the medical field it has become possible to analyze very small amounts of biological substances, and this contributes to early stage diagnosis for various diseases. For instance, malignant tumors can be diagnosed at an early stage by measuring alpha-fetoprotein and carcinoembryonic antigen. Further, immunological diseases can be diagnosed by measuring abnormal secretion of hormones such as insulin and thyroxine and immunoglobulin. Moreover, it has also become possible to effect monitoring during the course of medical treatment. Further analysis for low molecular hapten (incomplete antigen) such as medicines, contributes to plan a dose schedule of medicines.\nMost of these biological substances have been analyzed by an immunological method utilizing the antigen-antibody reaction. There have been proposed various immunological analyzing methods. For instance, antigen or antibody is fixed on fine particles made of glass or synthetic resin and is reacted with a sample antibody or antigen, while labeling antibody or antigen labeled with a highly sensitive marker such as radio isotope, fluorescent substance, luminescent substance and enzyme is used to produce antigen-antibody complex. Then the antigen-antibody complex is detected to quantify the sample antigen or antibody.\nFIG. 1 shows a reaction process of known enzyme immunological analysis utilizing an enzyme labeling reagent. On an outer surface of an insoluble carrier 1 there is fixed antibody 1A which is specifically reacted with antigen contained in a sample to be analyzed. If sample antibody instead of antigen is to be measured, antigen which is specifically reacted with the sample antibody has to be fixed onto the insoluble carrier 1. At first, the antibody 1A on the carrier 1 is reacted with antigen 2 contained in the sample. Then free antigen which has not been bound with the antibody 1A on the carrier 1 is separated from antigen 1A bound on the carrier 1 by means of washing. This separating process is generally called B-F separation. After that, the carrier 1 is reacted with a labeling reagent 3 comprising antibody 3A labeled with enzyme 3B, said antibody 3A being specifically reacted with the sample antigen 2. Next, the B-F separation is carried out again to separate free labeling reagent from labeling reagent bound with the sample antigen 2 reacted with the antibody 1A fixed on the carrier 1. Further, a color reagent is added to effect a color reaction under the existence of the enzyme 3B of the labeling reagent 3. Finally the reaction liquid is colorimetered to detect an enzyme activity. In this manner the sample substance 2 is quantified at a high sensitivity. The above mentioned immunological analyzing method has become widely practiced.\nIn the known immunological analyzing method using a labeling reagent containing antibody or antigen, after the antibody or antigen fixed on the carrier has been reacted with the sample antigen or antibody, the first B-F separation is performed, and then after the sample antigen or antibody has been reacted with the labeling reagent, the second B-F separation is carried out. Therefore, the analysis requires a larger number of steps and the analyzing process is liable to be complicated. Moreover, it is necessary to use various kinds of reagents and thus the running cost becomes expensive. It should be noted that an analyzer for effecting the known analysis becomes also complicated and expensive.\nFurther, since the substance to be analyzed is detected indirectly by using the labeling reagent as intermediate substance, the analysis is liable to be affected by slight change in external condition during the analyzing process. For instance, if the washing process is not carried out correctly, there might be produced measuring errors.\nThe immunological analysis has been also used to identify or determine a blood type. For instance, in German Patent No. 3,246,873, there is disclosed a method of determining the blood type by photoelectrically detecting a particle agglutination pattern. In this known method, blood cells are first separated by centrifuge, and then separated blood cells are dispersed in a saline solution to form a blood cell solution having a concentration of 2 to 5%. The blood cell solution is delivered into a reaction vessel having a conical bottom and a given antiserum is delivered into the reaction vessel. A particle pattern formed on the conical bottom is uniformly illuminated and is photoelectrically detected. A photoelectric output signal is processed to judge agglutinated or non-agglutinated particle pattern to determine the blood type.\nIn the above mentioned known blood type determining method, there are required various processes such as the centrifugal separation of blood cells, preparation of the blood cell suspension and the delivery of the antiserum. Therefore, the analyzing process is very cumbersome. Moreover, in order to form the particle pattern the reaction vessel has to remain still for a long time period. Further, in order to detect the particle pattern, there must be arranged correctly a light source, imaging lens, a plurality of light receiving elements, etc. Thus, the whole analyzing apparatus is liable to be large in size, complicated in construction and expensive in cost."}
{"text": "1. Field of the Disclosure\nThe subject disclosure relates generally to oilfield drilling, and more particularly to bottom hole assemblies and tools for orienting a bottom hole assembly (BHA).\n2. Background of the Related Art\nIn conventional drilling, the BHA is lowered into the wellbore using jointed drill pipes or coiled tubing. Often the BHA includes a mud motor, directional drilling and measuring equipment, measurements-while-drilling tools, logging-while-drilling tools and other specialized devices. A simple BHA having a drill bit, various crossovers, and drill collars is relatively inexpensive, costing a few hundred thousand US dollars, while a complex BHA costs ten times or more than that amount.\nMany drilling operations require directional control so as to position the well along a particular trajectory into a formation. Directional control, also referred to as “directional drilling,” is accomplished using special BHA configurations, instruments to measure the path of the wellbore in three-dimensional space, data links to communicate measurements taken downhole to the surface, mud motors, and special BHA components and drill bits. The directional driller can use drilling parameters such as weight-on-bit and rotary speed to deflect the bit away from the axis of the existing wellbore. In some cases, e.g. when drilling into steeply dipping formations or when experiencing an unpredictable deviation in conventional drilling operations, directional-drilling techniques may be employed to ensure that the hole is drilled vertically.\nDirection control is most commonly accomplished through the use of a bend near the bit in a downhole steerable mud motor. The bend points the bit in a direction different from the axis of the wellbore when the entire drill string is not rotating. By pumping mud through the mud motor, the bit rotates though the drill string itself does not, allowing the bit alone to drill in the direction to which it points. When a particular wellbore direction is achieved, the new direction may be maintained by then rotating the entire drill string, including the bent section, so that the drill bit does not drill in a direction away from the intended wellbore axis, but instead sweeps around, bringing its direction in line with the existing wellbore. As it is well known by those skilled in the art, a drill bit has a tendency to stray from its intended drilling direction, a phenomenon known as “drill bit walk”. A device for addressing drill bit walk is shown in U.S. Pat. No. 7,610,970 to Sihler et al. issued Nov. 3, 2009, which is incorporated herein by reference.\nThe use of coiled tubing with downhole mud motors to turn the drill bit to deepen a wellbore is another form of drilling, one which proceeds quickly compared to using a jointed pipe drilling rig. By using coiled tubing, the connection time required with rotary drilling is eliminated. Coiled tube drilling is economical in several applications, such as drilling narrow wells, working in areas where a small rig footprint is essential, or when reentering wells for work-over operations.\nIn coiled tubing drilling, a BHA with a mud motor is attached to the end of a coiled tubing string. Typically, the mud motor has a fixed or adjustable bend housing in order to drill deviated holes. Because the coiled tubing is unable to rotate from surface, a so called orienter tool is used as part of the BHA to “orient” the bend of the mud motor into the desired direction. There exists a multitude of different designs for the drive systems of such tools. Some designs support continuous rotation such as electric motor and gearbox drives, while others only permit rotation by a certain limited angle. The orienter tool is typically a high-torque, low-speed device, wherein the design of the drive system provides a torque output which can at least match the reactive torque exerted by the drilling mud motor.\nFor example, some orienter tools have utilized planetary gears in an effort to drive the output shaft. Basically, creating a torque on an output shaft means that a tangential force has to be exerted. By way of example, an output torque of 1,000 ft-lbs from a 2-inch diameter shaft means a tangential force of 12,000 lbs. This amount of force will quickly yield any material unless the tangential force is evenly distributed over a sufficient area to reduce the stress levels. In a conventional planetary stage with a size constraint on the order of 3 inches in diameter, the limits of how much bending force the gear teeth can take, and how much stress the planet carrier is capable of supporting will be much below 1000 ft-lbs of torque."}
{"text": "Portable electronic still cameras are widely used, which convert an optical image of a subject to an electric image signal through a CCD (charge-coupled device) image sensor or the like, convert the image signal to digital image data, and write the image data in a memory. The memory is capable of storing several tens of frames of image data, and it is possible to read out the image data therefrom for displaying or reproducing the subject image on a liquid crystal display (LCD) panel that is mounted to the camera body, or on external apparatuses such as a monitor CRT and a personal computer.\nAs the memory for recording the image data therein is generally used a DRAM (dynamic random access memory) which is accessible at high speed, and is usually called a flash memory. Because the flash memory is able to delete or rewrite the image data, it is easy to delete unnecessary image data from the memory or revise the memory with new image data. Using the recording medium repeatedly is a feature of the electronic still camera that is advantageous over those cameras using photographic film. In addition to the flash memory or other built-in memories like hard discs, it is possible to use a card memory or a flexible magnetic sheet that is removably attachable to the electronic still camera.\nMeanwhile, video printers have made it possible to make a hard copy of an image of a subject on the basis of image data written in the flash memory or the like. The video printers include Laser printer, thermal transfer printer, ink-jet printer and so forth. To make a hard copy, the image data read out from the flash memory of the electronic still camera is directly transferred to these printers, or is written in a recording medium like a card memory or a magnetic sheet, and then read out by the printers from the recording medium.\nSince the conventional video printers mainly use heat energy for printing, they consume certain amount of electric power. Besides, the printer body is so large that it is unhandy to carry about. Therefore, not being able to get a hard copy of a photographed image instantly is disadvantage to the conventional electronic still camera.\nIn view of the foregoing, an object of the present invention is to provide a portable electronic still camera having a printing device integrated therein that facilitates making a hard copy of a photographed image instantly.\nAnother object of the present invention is to provide an instant printer for printing an image on an instant film at a high speed on the basis of digital image data, that is so small and consume less electric power that it is handy to carry about or for integration into a camera.\nA further object of the present invention is to provide an instant film that is convenient for use in the instant printer of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a developing device used in an image forming apparatus such as an electrophotographic copier, a printer, a facsimile, and a composite apparatus having functions of those devices, and to an image forming apparatus provided with the developing device. In particular, the present invention relates to a developing device which replenishes a two-component developer including toner and carrier and discharges surplus developer and to an image forming apparatus provided with the developing device.\n2. Description of Related Art\nIn image forming apparatuses, an electrostatic latent image formed on an image carrier including a photosensitive member and the like is developed by a developing device and visualized as a toner image. Examples of the developing device include one employing a two-component developing method in which a two-component developer is used. The developing device of this type includes a developing container in which a two-component developer including toner and carrier is stored, and there are arranged a developing roller for supplying the developer to the image carrier and a stirring member for supplying the developer in the developing container to the developing roller while stirring and conveying the developer.\nIn the developing device, the toner is consumed by a developing operation; meanwhile, the carrier remains in the developing device without being consumed. Accordingly, the carrier stirred together with the toner in the developing container is deteriorated in proportion to a stirring frequency. As a result, charging performance of the carrier with respect to the toner is gradually deteriorated.\nIn this context, there has been well-known a developing device according to a first related art in which deterioration of charging performance is suppressed by replenishment of developer including carrier into a developing container and discharge of surplus developer.\nIn the developing device according to the first related art, two stirring members each including a rotary shaft and a helical blade helically formed about an outer periphery of the rotary shaft are arranged in parallel with each other in respective conveyance paths. A partition portion is provided between the conveyance paths, and communication portions for exchanging developer are provided to both end portions of the partition portion. A developer discharge port is provided on a downstream side of the conveyance path with respect to a developer conveying direction. Between the stirring member and the developer discharge port, a reverse helical blade helically formed in a direction reverse to that of the helical blade of each of the stirring members is provided as a discharge regulating portion integrally with the rotary shaft. With this structure, when being replenished into the developing container, developer is conveyed to the downstream side of the conveyance path while being stirred by rotation of the stirring members. When being rotated in the same direction as that of the stirring members, the reverse helical blade imparts a conveyance force in a direction reverse to the developer conveying direction due to the stirring members to the developer. The developer is retained by the conveyance force in the reverse direction on the downstream side of the conveyance path and increased in height. As a result, surplus developer climbs over the reverse helical blade (discharge regulating portion) so as to move to the developer discharge port, with the result of being discharged outside.\nHowever, in the developing device according to the first related art, even when new developer is not replenished, the developer conveyed by the helical blade of each of the stirring members moves to the downstream side of the conveyance path in an undulating manner in conformity with an outer periphery of the helical blade, with the result of colliding against the discharge regulating portion. When the developer collides against the discharge regulating portion, a height of the developer with respect to an outer periphery of the discharge regulating portion varies from each other in accordance with axial positions of the helical blade with respect to the discharge regulating portion. When the developer collides against the discharge regulating portion at a position at which the developer has height, the developer climbs over the discharge regulating portion so as to move into the developer discharge port by an impact of the collision. As a result, the developer is excessively discharged, which may lead to a risk of instability of the developer amount in the developing container. In particular, in an apparatus performing high-speed image formation, the stirring members are rotated at high speed together with the photosensitive member, and hence there prominently occurs an inconvenience of excessive discharge of the developer."}
{"text": "Natural gas is a widely available fuel which is primarily methane. The prices of natural gas has declined in recent years and are expected to possibly remain low for decades, natural gas has become a favored form of energy. As with any hydrocarbon based energy source, questions are raised about its contribution to greenhouse gases.\nGreenhouse gases are part of the ‘greenhouse effect’ that presents an environmental issue associated with the potential for global climate change. A number of gases in our atmosphere regulate the amount of heat retained close to the earth's surface. Water vapor, carbon dioxide, methane, and nitrogen oxides are the most common naturally occurring greenhouse gases. As the concentration of greenhouse gases increase, temperatures throughout the environment will rise significantly.\nThe continually increasing use of fossil fuels has caused an increase over the naturally occurring amounts of these gases and ways are sought to reduce the man-made influence in the generation of such gases. Carbon dioxide is seen as a major source of greenhouse gas. Even though carbon dioxide has less of an effect on trapping heat than the other naturally occurring gases, its prevalence from the increased reliance on hydrocarbons makes it suspect as a significant contributor to trapping of heat in global warming. For this reason those concerned with climate change are especially interested in limiting the generation of carbon dioxide.\nFrom the standpoint of CO2 production from hydrocarbon usage, natural gas offers important environmental benefits relative to the combustion of other fossil fuels. In rough numbers combustion of treated and unconverted natural gas will produce 30% less CO2 emissions than oil and almost 45% less CO2 emissions than coal. However, methane in the form of an emission is over 20 fold more potent in its greenhouse gas effect than carbon dioxide.\nWhile there are many applications for using natural gas as an energy source in its gas phase, attention has focused on how to utilize this energy source in the most efficient manner. Major drawbacks to the use of natural gas in traditional applications are problems of its storage and transportation as long as it remains in a gas form.\nOne way in which those seeking increased the natural gas utilization have attempted to overcome these problems is by employing processes for converting natural gas to denser and more easily transportable forms. Natural gas is typically converted into a liquid form by gas to liquid (GTL) processes, the primary of which either produces syncrude that directly converts the gas or produces synthesis gas that indirectly converts the natural gas. Although the combustion of natural gas in the unconverted gaseous form has green-house gas advantages, further processing of natural gas to liquids to transform it to more easily transported forms increases its greenhouse gas impact.\nOne solution to the problem of CO2 production from industrial processes is to sequester CO2 that is produced by such processes. Effecting carbon dioxide (CO2) capture and sequestration (CCS) has been proposed by various technologies that can reduce CO2 emissions from facilities such as coal-fired and gas-fired power plants, and large industrial sources. Most CCS technologies employ three steps. The first step will capture CO2 from oil refining operations, treatment of well head gas, powerplants, renewable fuel production or other industrial processes. Compression of the captured CO2 prepares it for transfer, in tanks, pipelines or via other means, to permanent sequestration facilities. These facilities usually comprise geologic formation located a mile more underground that hold the CO2 in a porous layer of rock situated beneath a gas impermeable layer of rock. The impermeable layer of rock prevents an upward migration of the CO2. This type of CO2 sequestration is expensive, used infrequently and impractical for major sources of CO2 emissions.\nAnother major source of CO2 emissions is in the use of hydrocarbons as motor fuels. An approach to reduce the effect of CO2 emissions from hydrocarbon fuel combustion is to produce motor fuels from renewable resources such as plant materials. This approach reduces the use of non-renewable, i.e. fossil, hydrocarbons that when combusted put carbon dioxide into the atmosphere by releasing the fossil hydrocarbons from their naturally sequestered state. In contrast, using renewable sources for feedstocks, i.e. plant materials, maintains a closed CO2 system where the carbon dioxide released by hydrocarbon combustion balances with the carbon dioxide captured by the plant material that goes into the feedstock for producing the fuels and chemicals.\nPlant materials for use in such feedstocks can come from almost any plant source. One general type of plant material, such as corn, rice, cane, wheat, sorgum etc., have high carbohydrate concentrations that are readily converted to sugars. At this time essentially all of the production of ethanol for motor fuel purposes relies on the saccharization of plant materials into sugar and its fermentation to alcohols. These fermentations produce CO2 directly as a by-product of the fermentation process. In such processes the amount of CO2 produced directly can equal the mass of ethanol. Moreover, concerns about plant sources can arise since their use requires utilization of land and water resources for growth of these corps or trees and is subject to limitations posed by alternative uses of these land and water resources for food and other agricultural and forest product production.\nAnother way to produce alcohol motor fuels or other chemicals is by generating syngas from plant material or natural gas. In the case of plant material, its gasification will produce a syngas. In the case of natural gas it may be converted to reformed gas. The syngas obtained by gasification or by conversion of natural gas are suitable substrates for anaerobic fermentations of hydrogen and carbon monoxide. Anaerobic fermentation processes involve the contact of the substrate gas in an aqueous fermentation menstruum with microorganisms capable of generating alcohols such as ethanol, propanol, i-butanol and n-butanol. The production of these alcohols requires significant amounts of carbon monoxide and carbon dioxide together with hydrogen.\nIn addition to the aforementioned plant and natural gas sources, a suitable source of the substrate gas for carbon monoxide and hydrogen conversions can be derived from the gasification of carbonaceous materials other than those specifically mentioned such as the partial oxidation of natural gas, biogas from anaerobic digestion or landfill gas, coal coking and industrial steel manufacture. The substrate gas containing, carbon monoxide, hydrogen, and carbon dioxide, usually contains other components such as water vapor, nitrogen, methane, ammonia, hydrogen sulfide and the like. (For purposes herein, all gas compositions are reported on a dry basis unless otherwise stated or clear from the context.)\nIn any of the above cases certain anaerobic bacteria, especially those from the genus Clostridium can produce various alcohols, particularly ethanol, from the resulting substrates comprising CO, CO2 and H2 via the acetyl CoA biochemical pathway. Different strains of Clostridium ljungdahlii that exemplify producing ethanol from gases are described in WO 00/68407, EP 117309, U.S. Pat. Nos. 5,173,429, 5,593,886, and 6,368,819, WO 98/00558 and WO 02/08438. Abrini et al, Archives of Microbiology 161, pp 345-351 (1994)\nThe effectiveness of converting CO2 to other products will depend on the amount of CO and H2 that is present with the CO2 in the converted gas stream. For instance, the theoretical equations for the conversion of carbon monoxide and hydrogen to ethanol are:6CO+3H2O→C2H5OH+4CO2 6H2+2CO2→C2H5OH+3H2O.As can be seen, the conversion of carbon monoxide results in the generation of carbon dioxide. The conversion of hydrogen involves the consumption of hydrogen and carbon dioxide, and this conversion is sometimes referred to as the H2/CO2 conversion. The microbiological processes that take place have been expressed for convenience in the formulas shown above. It is known that in the microbiological process the production of the alcohol requires a carbon source and a biological energy source. The formulas show that the biological energy source is the available electrons in the hydrogen and carbon monoxide of the substrate gases. As a result gas streams with high CO will produce more CO2 or those with low H2 will limit the conversion of CO2 and in either case may make the anaerobic process a net producer of CO2.\nRegardless of originating source, the substrate gases are typically more expensive than equivalent heat content amounts of the carbonaceous feedstock fuel. Hence, a desire exists to use these gases efficiently to make higher value products. The financial viability of any conversion process, especially to commodity chemicals such as ethanol, will depend, in part, upon the costs of the feedstocks, conversion efficiency and operating and capital costs for generating the substrate gases; and upon the capital costs, the efficiency of conversion of the carbon dioxide, carbon monoxide and hydrogen to the sought products, the process energy costs and the extent of the electron availability and utilization to effect the conversion of the substrate gases to the higher value products.\nIn a bioreactor, hydrogen and carbon oxides pass from the gas phase to being dissolved in the aqueous menstruum, and then the dissolved hydrogen and carbon oxides contact the microorganisms for bioconversion. Due to the low solubilities of carbon monoxide and, especially, hydrogen in aqueous media, mass transfer can be a factor limiting rate and conversion in the bioconversion to alcohol.\nThe off gases from bioreactors contain substrate that was not bioconverted and diluents such as methane and nitrogen. Although off gases can be recycled to the bioreactor or passed to another bioreactor, challenges can exist. For instance, the substrate gases may contain diluents that, if recycled to a bioreactor, can build-up and reduce the partial pressure, thus reducing the driving forces for mass transfer of hydrogen and carbon monoxide to the aqueous menstruum.\nOverall, the net effect of the processing as gas into liquids such as alcohol or hydrocarbons will produce an increase in the CO2 emissions when compared to a liquid to liquid conversion such as petroleum crude to gasoline. No one to date has disclosed a process that can improve the overall greenhouse gas emissions from the processing of natural gas to produce ethanol.\nBell in United States Published Patent Application No. 20100105118 discloses an integrated process for making alcohols from natural gas which provide a high bioconversions of carbon monoxide in fermentations in the absence of oxygen. Bell notes at paragraph 0013 that, in theory, carbon dioxide may be used as a reactant for the production of higher alcohols such as ethanol. However, he states that in practice the fermentation route to higher alcohols tends to be a net producer of carbon dioxide. In his disclosed process, the gas from the bioreactor which contains carbon dioxide is fed to a reaction section of a steam reformer. The reformer is either operated dry or with a mole ratio of water to carbon dioxide of less than 5:1. Bell confirms the low conversion of hydrogen in the examples. Although Bell may have reduced carbon dioxide emissions as compared to the use of autothermal reforming or traditional steam reforming, the low conversion of hydrogen detracts from the commercial viability of the disclosed process.\nProcesses are therefore sought that can reduce the carbon footprint in the utilization of the non-renewable resources.\nProcesses are also sought that can reduce emitted greenhouse gases by utilizing CO2 released in the fermentation of renewable resources.\nProcesses are also sought that can maximize the production and availability of electrons for microorganism growth and alcohol production while also reducing CO2 emissions for the production of alcohol in commercial-scale, continuous operations."}
{"text": "1. Field of the Invention\nThe invention relates to an apparatus for temporarily increasing the load-bearing capacity of bridge girders, of the kind used in bridge cranes or polar bridge cranes.\nWhen heavy components are lifted, an installed load-bearing capacity of existing bridge girders is often less than a capacity which is actually required. Such a situation arises, for instance, in unusual conversion or repair provisions in a power house of a power plant or other industrial plants, for instance when generators or steam generators in nuclear power plants are being replaced. Until now, in order to lift heavy machine parts or when placing other unplanned stresses on bridge girders, relief devices such as central support masts that are each constructed for a particular application which is involved have been put in place, or else separate hoisting tools have been installed at high cost.\n2.Summary of the Invention\nIt is accordingly an object of the invention to provide an apparatus for temporarily increasing the load-bearing capacity of bridge girders, which overcomes the hereinafore-mentioned disadvantages of the heretofore-known devices of this general type, which is universally usable and which is easily mounted.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, an apparatus for temporarily increasing a load-bearing capacity of a bridge girder, comprising portable supports having deflection pulleys; a relief cable associated with the pulleys; the supports, the pulleys and the cable temporarily mounted above bracing points of a bridge girder on both sides of a support point; and a load-bearing element carried by the relief cable for relieving the bridge girder at the support point.\nTherefore, the object of temporarily increasing the load-bearing capacity of crane bridges is attained according to the invention by using a portable cable structure to reinforce the bridge girder.\nIn accordance with another feature of the invention, the relief cable has two ends and tensioning devices securing the two ends to the bracing points or to the bridge girder.\nThe supports are mounted on both ends of the bridge girder, above the bracing points and protrude upward beyond the bridge girder. The deflection pulleys that serve to receive the relief cable are associated with upper ends of the supports.\nThe supports are connected to one another by the relief cable, which passes over the two deflection pulleys. The relief cable is connected to the bridge girder at least at one point between the two supports, preferably in the middle of the girder. This is carried out, for instance, through the use of a load-bearing element having a lower end, for instance, which reaches under the bridge girder whose capacity is to be increased and supports it from below at this point. Thus an upwardly-directed vertical force, which engages the bridge girder at the support point, is generated from the height of the associated supports and the cable tension of the relief cable, which is tightened with the aid of the tensioning devices, such as hydraulic cylinders, or screw elements. The forces occurring as the relief cable is tightened can be favorably introduced advantageously and statically at end surfaces of the bridge girder through the use of an abutment that engages the bridge girder from below.\nThe load-bearing element associated with the relief cable acts like a center support that supports the bridge girder from below with a vertical force.\nIn accordance with a concomitant feature of the invention, the bridge girder is disposed in a power house of a power plant or in a reactor building of a nuclear power plant.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in an apparatus for temporarily increasing the load-bearing capacity of bridge girders, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawings."}
{"text": "This invention is in the field of haptic devices. Embodiments of the invention are more specifically directed to remote force management by operation of one or more haptic devices.\nConventional computer systems typically interface with their users, and with the “outside world” in general, by way of graphics and text. Even pointing devices (e.g., mouse, trackball, trackpad) operate in a graphical sense, by navigating a cursor among graphical elements. In the industrial context, additional interface approaches include the receiving of inputs from sensors and transducers, and the outputting of motor control signals in either of the digital and analog domains.\nHaptic devices refer to those interface devices by way of which computers and computer systems mechanically interface with users or objects, by applying forces, vibrations, or motions from the device to the user or object. Examples of conventional haptic devices include robotic devices for handling or moving an object, for example as controlled by a haptic input device receiving motion inputs from a human user, vibrators in game controllers and other devices for issuing a sensation to the user at the appropriate times in a game or simulator, virtual reality systems by way of which user motion is directly interpreted as inputs, and by way of which the haptic device can receive physical feedback, and the like.\nConventional haptic devices operate according to a variety of mechanisms. A common type of haptic device includes actuators that cause mechanical motion in response to a signal stimulus. Conventional haptic input mechanisms include electroactive polymers for sensing the touch of a user, and actuators including piezoelectric, electrostatic, and subsonic audio wave elements for transducing touch or pressure to an electrical signal.\nIt is contemplated that multiple conventional haptic devices may operate in conjunction with one another by way of a shared motor mechanism, such as a lever. However, this shared motor mechanism necessarily enforces a practical limit as to how remote the devices can be from one another, and as to the response of the devices to one another."}
{"text": "The present invention relates to a method of controlling exhaust-gas recirculation (EGR) particularly for an internal combustion engine having primary and secondary intake passages.\nThere have been known internal combustion engines which include a primary intake passage for supplying an air-fuel mixture under a full range of loads and a secondary intake passage for supplying an air-fuel mixture under high loads. The primary intake passage has a diameter that is small enough for it to be able to take care of the practical range of engine operation, and the secondary intake passage has a larger diameter such that it will come into operation only when the engine undergoes acceleration or is otherwise put under a higher load. Such a larger ratio between the cross-sectional areas of the primary and secondary intake passages is selected for more effective functioning of the internal combustion engine with primary and secondary intake passages. With the larger cross-sectional ratios of the primary and secondary passages, however, an EGR valve for the engine cannot be effectively actuated only by a vacuum introduced through a vacuum signal port from the primary intake passage, and hence a sufficient amount of EGR which is proportional to engine loads is not available."}
{"text": "The present invention generally relates to a game field and method for playing competitive paintball. More specifically, the present invention relates to a method and set of rules for playing competitive paintball that increases the amount of game action and scoring to create a game that appeals to spectators.\nThe game of paintball generally includes two separate teams of players that each have paintball guns. The paintball guns use compressed gas to propel paint-filled gelatin or plastic shells at players from the opposing team. When a paintball strikes an opposing player, the shell breaks and leaves a splatter mark of paint on the player, indicating that the player has been eliminated.\nCurrently, paintball competitions are played using one of two basic formats --\"center flag\" or \"dual flag\". In a center flag competition, a single flag is hung at the middle of the playing field. Two competing teams start from different ends of the playing field and attempt to capture the flag and carry it to their opponent's flag station on the opposite end of the field. In the dual flag system, each of the opposing teams begins the game at opposite ends of the field with their team's flag hanging at their own flag station. The object of the game for each team is to capture the opposing team's flag and return it to their team's station while at the same time protecting their team's flag from capture from the opposing team. In each of these games, a team player is eliminated by being shot by a paintball fired from a member of the opposing team.\nIn the conventional paintball games described above, points are typically awarded for the number of opponents eliminated during the game and for the first team to grab the opponent's flag, with additional points being given to the first team to hang their opponent's flag in the proper station. Each paintball match is played until one of the teams captures their opponent's flag and returns it to the proper flag station, or until the game time expires. Since points are awarded to both teams based on the number of opponents eliminated, a premium is placed on avoiding elimination. As an illustrative example, if a player gets eliminated, the player is giving points to the opposing team, regardless of which team is successful in capturing the flag. This type of scoring places a burden on the individual players to avoid being eliminated by making risky moves or engaging in risky strategies of play. Therefore, paintball games typically evolve into a series of long-range shooting duels, frequently with little or no action for long periods of time.\nBecause players on each of the opposing teams in a conventional paintball competition tend to remain behind stationary bunkers or obstacles and engage in long-range shoot-outs, paintball competitions have not been able to attract a large number of spectators, either for live action events or televised broadcasts of competitive tournaments. Therefore, a need exists for a paintball game including a set of rules that encourages a fast-paced, high-scoring competition that allows viewers to follow the action and be entertained by risky and daring playing strategies.\nIt is an object of the present invention to provide a method of playing a paintball competition that rewards players for aggressive and risky movements between bunkers organized on a playing surface. It is an additional object of the invention to provide a scoring system that rewards the scoring team for scoring a goal, while determining the number of points awarded based upon the position of the opponent's flag when the scoring team scores. It is an additional object of the present invention to provide a paintball game that includes various penalties for improper game play and removes a player from the game for committing one of specified penalties. It is an additional object of the present invention to place a time limit on the paintball game that is divided into several individual timed periods. Further, it is an object of the invention to provide a paintball game that returns each team to full strength after each score and upon completion of each period, thereby encouraging the players from each team to make daring moves and score as many points as possible during each period. Finally, it is an object of the present invention to provide a paintball game having a set time limit as opposed to prior games having a varying time limit."}
{"text": "This invention relates to an atomic force microscope for profiling the properties of a surface at nanometer resolution and for probing the properties of individual molecules attached to that surface.\nThe atomic force microscope utilizes a sharp probe which may be precisely positioned above a surface. The probe is attached to a positioning device via a soft (spring-like) cantilever. Interaction with the surface is indicated by motion of the probe end of the tip with respect to the driven end of the cantilever. In its common use for mapping the topography of a surface, this deflection is held constant by adjusting the probe to sample distance as the probe is scanned over the surface.\nThe interaction between the probe and the surface is of considerable interest in its own right. By sweeping the tip back and forth from the surface and monitoring the deflection of the probe relative to the driven end of the cantilever, so called force curves may be mapped out as taught by Elings et al., U.S. Pat. Nos. 5,224,376; 5,237,859; and 5,329,808. These curves are valuable both for setting the operating point of the microscope and for mapping the elastic properties of a surface.\nThe scheme for obtaining force curves according to the prior art is set out in FIG. 1A. A piezoelectric actuator 1 is used to sweep the stiff end of a cantilever 2 by some amount (labeled X.sub.D) toward the sample S in FIG. 1A. The flexible cantilever 3 (having a stiffness k.sub.T) bends as the tip 4 indents the sample, represented in FIG. 1A by a spring 5 of stiffness k.sub.s. As a consequence, the displacement of the tip, X.sub.T, is generally less than the displacement of the stiff end of the cantilever, X.sub.D. This situation is represented in FIG. 1B which shows the arrangement of the two springs corresponding to the cantilever (k.sub.T) and the sample (k.sub.S) together with the displacement of the actuator (X.sub.D) and the tip (X.sub.T). In equilibrium, the force, F, along the mechanical path must be equal at all points and, far below resonant frequencies where the cantilever and sample behave as Hookean springs: ##EQU1## Here k.sub.eff is an effective spring constant describing the relationship between displacement and applied force at the point at which the force is applied. From Equation (2) we see that ##EQU2## Thus, in this method for measuring the stiffness, the effective spring constant is equivalent to a series combination of the two springs represented by k.sub.S and k.sub.T. In such a relationship, the weakest spring dominates the spring constant of the combination. Thus, this prior art scheme is of limited use for measuring the stiffness of soft samples, unless a correspondingly soft (i.e., having greater flexibility) cantilever is used. The use of soft cantilevers is often precluded by their increased propensity to jump into contact with the surface.\nTo make matters worse, the prior art method relies on sweeping either the entire sample up towards the tip, or the entire tip assembly (and scanning tube) down towards the sample. These massive components have a low resonant frequency by virtue of their mass, so this method is intrinsically slow, a disadvantage when it is desirable to obtain force curves rapidly.\nA further disadvantage of this method of obtaining force curves lies with the piezoelectric material that is used as the displacement actuator. This material suffers from a considerable hysteresis. On reversing the direction of the applied voltage, the cantilever displacement usually lags the displacement at the same applied voltage but reached by scanning from the opposite direction. This is illustrated in FIG. 2 which shows a typical force curve. This is a plot of cantilever displacement versus distance moved towards or away from the surface. When the tip is far from the surface, the cantilever is not deflected so the plot is a flat horizontal line 6. When the tip is close enough to the surface, it often snaps uncontrollably into the surface (depicted at 7), remaining in contact and being pushed up by the surface (line 8). When the direction of the sweep is reversed, the tip comes down with the surface (line 9), but the scan is displaced by some amount .DELTA.z along the horizontal axis owing to piezoelectric hysteresis. The cantilever tip pulls off the surface again (depicted at 10), further displaced from the original contact point by adhesion between the tip and surface. The amount of the hysteresis depends upon the sweep rate and other characteristics of the piezoelectric actuator and its history. It is a poorly characterized quantity and the source of considerable uncertainty in obtaining force curves.\nAnother method for driving the tip has been described by Lindsay et al., \"Scanning Tunneling Microscopy and Atomic Force Microscopy Studies of Biomaterials at a Liquid-Solid Interface,\" J. Vac. Sci. Technol. 11:808-815 (1993); Lindsay, U.S. Pat. No. 5,513,518; and Lindsay, U.S. Pat. No. 5,515,719; and also by Florin et al., \"Atomic Force Microscope with Magnetic Force Modulation,\" Review of Scientific Instruments 65:639-643 (1993) and O'Shea et al., \"Atomic Force Microscopy of Local Compliance at Solid-Liquid Interfaces,\" Chemical Physics Letters 223:336-340 (1994). In this approach, a force is applied directly to the tip (as opposed to the rigid part of the cantilever holder) by fixing a magnetic particle or film on the tip and using an external magnetic field to drive the tip. This is illustrated schematically in FIG. 3A. Here, the actuator 1 is held in a fixed position, as is the rigid part of the cantilever assembly 2. The magnetic particle or film 12 is acted on by an external magnetic field so to exert a force F on the end of the flexible cantilever 3, pushing the tip 4 into the sample represented by the spring 11 of stiffness k.sub.S. In this case, referring to FIG. 3B, both the stiff end of the cantilever and the sample are fixed so EQU F=(k.sub.S +k.sub.T)X.sub.T =k.sub.ef X.sub.T (4)\nwhere EQU k.sub.eff =k.sub.T +k.sub.S. (5)\nSo, in this case, the cantilever and sample act as springs in parallel. This method has the advantage that a soft surface (k.sub.S small) does not lower the resonant frequency of the assembly appreciably. Furthermore, because only the tip is being moved, and this has by far the smallest mass of all the microscope components, rapid operation is possible.\nThere are two methods for driving the tip magnetically. In one (FIG. 4A) a particle of a magnetic material 12 is placed on the end of the cantilever 13 with its magnetic moment M aligned along the direction of desired motion. A magnetic field gradient, dB/dz, 14 is applied along the same direction as the magnetic moment, resulting in a force, F.sub.z along the z direction ##EQU3## A second geometry (FIG. 4B) utilizes a film having its magnetic moment M (15) perpendicular to the direction of desired motion. A magnetic field B is applied in the direction of desired motion 16 so that a torque, n is generated on the cantilever according to EQU n=M.times. (7)\nand, in the geometry shown, this is equivalent to a force on the end of the cantilever in the direction of desired motion F.sub.z given by ##EQU4## where l is the length of the cantilever. This latter method of operation has been described by Lindsay, U.S. Pat. No. 5,515,718 and verified by Han et al., \"A Magnetically-Driven Oscillating Probe Microscope for Operation in Liquids,\" Applied Physics Letters 69:4111-4114 (1996). It is more sensitive than the former method which relies on the magnitude of dB/dz. Large magnetic field gradients are difficult to generate. In the second method, where a given B field produces a torque, this is translated into a large force by the shortness of the cantilever, as described by Equation 8.\nThis method is most sensitive when operated in an AC mode as taught by Lindsay and Han et al, supra, and as demonstrated by the force curves obtained by O'Shea et al., supra. This method of operation is outlined in FIGS. 5A and 5B. When the cantilever 20 is far from the surface 21 (FIG. 5A), the tip is oscillated with an amplitude X.sub.0. If the driving frequency is well below a resonance frequency of the tip, then ##EQU5## whereas, when the cantilever 20 is interacting with surface 21 (FIG. 5B) of stiffness k.sub.S , ##EQU6## As the surface is approached, k.sub.S becomes much bigger than k.sub.T so that the ratio X/X.sub.0 becomes 1/k.sub.S or equal to the compliance of the sample or surface. It should be noted that a relationship similar to Equation 11 can be derived for the ratio X.sub.T /X.sub.D for the case where the stiff end of the cantilever is driven. However, this is not experimentally useful because it is the quantity X.sub.T that is measured (referred to as X in equation 11). Where X.sub.T is far from the surface, it is not simply related to X.sub.D if the tip is driven in a medium such as a fluid, even if far from resonance.\nFor the case of direct-force driving of the tip near a surface, a plot of the quantity X/X.sub.0 as a function of distance from the surface reflects the relative compliance of the surface, (k.sub.S /k.sub.T).sup.-1. Such a curve might more appropriately be called a compliance curve to differentiate it from a force curve. The inherently greater sensitivity and information content of compliance curves is illustrated by the measurements of O'Shea et al., supra.\nFinally, there is presently considerable interest in measuring the stiffness of individual molecules attached between the tip and the surface as illustrated in FIG. 6. Chemical means are used to attach one end of a macromolecule 30 to the tip 4, while the other end is attached to the surface 21 below the tip 4. The prior art has consisted of generating a force curve (like that shown in FIG. 2) as the molecule is stretched or compressed by motion of the piezoelectric actuator 1. An example of such a measurement is the recent work by Reif et al., \"Reversible Unfolding of Individual Titin Immunoglobulin Domains by AFM,\" Science 276:1109-1112 (1997). Dramatic variations in stiffness were observed as particular domains of a protein were unfolded sequentially.\nAccordingly, the need still exists in this art to provide an apparatus for measuring the compliance of samples on a nanometer scale as a function of the approach distance of a tip more rapidly and without the hysteresis problems of the prior art. There is also a need to provide a method for setting the operating point of an imaging microscope to a fixed value of surface compliance prior to scanning the tip over the surface."}
{"text": "The invention relates to an apparatus and a method for transmission of data over a network, and from point to point. More particularly, it relates to an apparatus and method for configuring a multi-purpose optical transmitter or transceiver for communication of video, audio, and data signals, each with differing protocols, on a fiber optic network capable of utilizing a single optical fiber to communicate any of the differing types of signals and protocols associated with it. Further, a software program is provided to allow users to configure or customize the system using menu-driven graphical interface screens and common configurable components.\nWith the ever-widening reach of public and private data networks used by the public, and especially by business and government, there continues to be an increase in the flow of electronic data over such networks. Such communications may be from remotely-located differing branches of government offices, or different divisions of a company. Further, such electronic communications may include video, audio or other types of data signals, each with its own communications protocols that must be followed in the transmission and reception of intelligible communications signals.\nAdditionally, with this rise in electronic traffic between remote locations worldwide, there is an ever-increasing demand for a transmission system that can both transport and accommodate increasing levels of bandwidth consisting of multiple channels of video, audio and/or data over long distances. The optimum solutions for this demand are systems which involve the least number of transmission mediums as possible and that are unaffected by environmental hazards from static and electrical disturbances.\nSatellite communications, such as those involving the frequency band of 390 MHz to 2150 MHz (often called the L-band) provide significant bandwidth. Once the L-band signal is received, it is increasingly desirable to use a fiber optic communication network (which provides sufficient bandwidth to accommodate data through-put of the order produced by satellite communications) to transmit that data terrestrially.\nIn general, fiber optics serve as the best solution to the problem of environmental interference. Transmission of electronic signals for data of all types over fiber optic networks will eliminate EMI radiation, EMI susceptibility and other factors which would otherwise affect copper or other cable. Consequently, this makes fiber the choice for high quality video transmission as well as conventional data transmission.\nHowever, as noted, there are a variety of video, audio and data signals, using protocols specific to the type of signal carried, that a multipurpose fiber optic transmission system may accommodate. Conventionally, for each of these signal types and their respective protocols, a generally unique transmission component is employed to generate the corresponding signal.\nBecause each user at each transmission and reception site may be employing any or all of the different electronic signals and protocols, it is a complex task to design and manufacture site-specific data transmission systems that will communicate with remote sites and which may be required to accommodate the same or other types of electronic data signals and associated protocols. Such differing types of data stream transmission and reception requirements (even between two remote locations of the same company or government department) can create logistical difficulties. Furthermore, designing a data transmission and reception system which will suit each customer's individualized needs can be complex and time consuming. In the case of satellite-derived transmissions, it may be desirable to transmit the signal across a fiberoptic communications network without changing the corresponding transmission protocol or otherwise changing the signal, together with video, audio or other data.\nCurrently, such designs may entail computers with multiple types of data transmission cards or components engaged to local and wide area networks. Generally, each type of data will require a corresponding card or component. Of course, each of the differing cards or components may also require software or drivers to allow it to communicate with the controller or computer with which it must send and receive signals according to the proper communications protocol. Further complications arise when signals originate over an Ethernet or satellite network, and must be converted to transmission over a fiber optic network for communication to the remote destination.\nTherefore, it can be seen that there is a pressing need for an apparatus, and a corresponding method, that provides a single system to communicate the various data types that might originate in any highly-customized client venue location, and that allows for transmission of those signals to one or more remote locations using a fiber optic network. Such a system should allow for any type of conventionally-employed data protocol to be used by the originator and receiver of the signals, yet still allow the signals to be transmitted over a fiber optic network.\nStill further, in a particularly preferred mode of the invention, a configuration of the manufactured system allows for customization by the buyer according to the existing communications streams and protocols being currently used by the buyer. The buyer is able to design a customized component, the component thereby being able to accommodate current communications data streams and protocols and to communicate them to remote locations using a single apparatus providing fiber optic transmission. Such customization, heretofore being hard to achieve if at all, allows users to maintain their current network communications types and protocols, while still being able to take advantage of the benefit of fiber optic network transmission. This eliminates any requirements for the users to upgrade and change all their systems and protocols before taking advantage of the benefit of fiber optic network communications, thereby eliminating the costs and inevitable problems that such an upgrade would entail.\nAs such, those skilled in the art will appreciate that the conception upon which this disclosure is based may readily be utilized as a basis for designing other methods and systems for carrying out the several purposes of the present invention. It is important therefore that the claims be regarded as including equivalent constructions insofar as they do not depart from the scope of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a downhole oil and gas well completion tool which is operatively connectible to a lower section of tubing or a packer and configured to perform multiple completion-related functions.\n2. Background Art\nThe harvesting of hydrocarbons from a subterranean formation involves the deployment of a drilling tool into the earth. The drilling tool is driven into the earth from a drilling rig to create a wellbore through which hydrocarbons are passed. Once a predetermined well depth is reached, the formation is tested to evaluate and determine whether the well will be completed for production, or plugged and abandoned.\nCompletion of a well generally refers to the operations that prepare a well bore for producing oil or gas from the reservoir. The goal of these operations is to optimize the flow of the reservoir fluids into the well bore, up through the producing string, and into the surface collection system.\nThe well bore is typically lined (cased) with steel pipe, and the annulus between well bore and casing is filled with cement. Properly designed and cemented casing prevents collapse of the well bore and protects fresh-water aquifers above the oil and gas reservoirs from becoming contaminated with oil and gas and the oil reservoir brine. Similarly, the oil and gas reservoir is prevented from becoming invaded by extraneous water from aquifers that were penetrated above or below the productive reservoir.\nThe nature of the reservoir, evaluated from a core analysis, cuttings, or logs, or from experience with similar productive formations, determines the type of completion to be used. In a barefoot completion, the casing is set just above the producing formation, and the latter is drilled out and produced with no pipe set across it. Such a completion can be used for hard rock formations which are not friable and will not slough, and when there are no opportunities for producing from another, lower reservoir. Set-through and perforated completions are also employed for relatively well-consolidated formations from which the potential for sand production is small. However, the perforated completion is used when a long producing interval must be prevented from collapse, when multiple intervals are to be completed in the one borehole, or when intervening water sands within the oil-producing interval are to be shut off and the oil-saturated intervals selectively perforated.\nA string of steel tubing is lowered into the casing string and serves as the conduit for the produced fluids. The tubing may be hung from the well-head or supported by a packer set above the producing zone. The packer is used when it is desirable to isolate the casing string from the produced fluids because of the latter's pressure, temperature, or corrosivity, or when such isolation may improve production characteristics. The string, which may be referred to herein as a tubing string, may comprise any number of components known in the art. Such components, in addition to tubulars, may include tools, joints, packers, etc.\nTo complete the well, casing is installed and cemented in the wellbore, then production tubing is installed in the casing, which is perforated so that hydrocarbons may pass from the formation into the wellbore, and up to the tubing string to the surface for collection. Often a series of tests are conducted as a part of this process, to confirm the integrity of the casing and tubing.\nWhen carrying out testing or other operations in a wellbore, test equipment or other apparatus may be mounted on an end portion of a string of tubular sections, known as tubulars to form a tubing string. The equipment is lowered into the bore on the end of the string, the length of the string being increased by the addition of further tubulars, which are threaded together to define a continuous internal bore between the apparatus and the surface.\nA number of traditional methods exist for testing completion of a well and tubing.\nThe most commonly used method involves the use of a wire-line retrievable plug. This process typically involves the hiring of a wire-line contractor to both run and pull the plug. The overall process is time-consuming, typically taking about 3-4 hours to set up the wire-line unit, and lower and set the plug by the wire-line. After setting the hydraulic packer and testing the completion and tubing, the wire-line operator will run the wire-line into the tubing again to retrieve the plug which might consume another 3-4 hours if no delays are encountered. Debris and impurities in the completion fluid and/or pressure trapped around the plug often result in sticking of the plug in the well. Retrieving a stuck plug can greatly increase the length of the process and may also lead to a loss of wire-line in the well, which will require the hiring of additional specialists to perform a fishing operation.\nThe use of the above-mentioned wire-line operation is typically feasible only if the well deviation is no greater than 70 degrees. If the testing location where the plug will need to be set is at a greater deviation, the wire-line method may not be practical and the operator must use another method of running the plug such as coiled tubing. A coiled tubing operation, once on location, takes about 5 hours to set up. Once all of the equipment is set up, running the plug by means of the coiled tubing can easily take at least 4 hours, depending on depth. The coiled tubing operation may easily cost tens of thousands of dollars in addition to the total time used to run and retrieve the plug which may exceed 10 hours, if no problems are encountered.\nA more recently developed method for setting a hydraulic packer and testing the completion and tubing involves the use of a glass disc which is run inside a special pipe attached to the bottom of the tubing string. When using the glass disc method, the tubing cannot be self filled, which will require that the completions operator either manually fill the tubing via a water hose from the surface, requiring significant time, or the operator will need to add another piece of equipment to the tubing string known as a self-filling tool.\nAfter setting the packer and testing the tubing the completions operator needs to break the glass disc in order to have the tubing opened and ready for production. If the well is not highly deviated, the wire line contractor can set up his equipment and run into the well with his tool string and break the glass disc. This operation will take around 4-5 hours.\nIf the well is highly deviated (more than 70 degrees) then the completions operator may need to use a coil tubing contractor to break the glass disc. Also this operation will take around 5 to 6 hours in addition to the coil tubing set up/rental cost.\nAnother method of setting the packer and testing the tubing is to use a pump out plug which is a special short pipe with a ball seat fixed in a seat with a number of shear screws. After running the tubing completely, a ball is dropped from the surface. It takes around 40-60 minutes until the ball seats on the ball seat, then surface pressure is applied against the ball to set the packer and test the tubing. After testing the tubing the surface pressure is increased to shear the ball seat shear screws and pump the seat and the ball down into the well bottom. This method holds pressure from above only and can not hold pressure from below. For this reason another plug/barrier is needed to be run at the bottom of the tubing while dismantling the rig blow out preventor (BOP) and mounting the Christmas tree. Also some problems can happen when shearing the shear screws of the pump out plug due to completion fluid pressure differential across the ball, that leads to inaccurate shear value (either more or less than the predetermined pressure for shearing).\nAccordingly, a need exists for a tool capable of performing various completion-related operations without requiring repeated pulling of the tool, or hiring of one or more specialists."}
{"text": "1. Field of the Invention\nThe invention relates to a lens cap, more particularly to a sliding lens cap for a lens unit of a projector.\n2. Description of the Related Art\nAs shown in FIG. 1, a conventional lens cap 10 for covering and uncovering a lens unit 101 of a projector 100 is connected by a cord 11 to a housing of the projector 100 so that the lens cap 10 can be prevented from being misplaced when the lens cap 10 is removed from the lens unit 101. When operating the projector 100, the lens cap 10 is removed from the lens unit 101 and is suspended from the housing by the cord 11. After using the projector 100, the lens cap 10 is replaced on the lens unit 101 to cover the same. To retain the lens cap 10 covering the lens unit 101, the lens cap 10 is provided with an engaging mechanism to engage removably a periphery of the lens unit 101 such that a force is required to disengage the lens cap 10 from the lens unit 101, which results in inconvenience. Moreover, if an operator forgot to remove the lens cap 10 before operating the projector 100, the lens cap 10 will be subjected to a strong light from the projector 100 that can melt the lens cap 10."}
{"text": "Anti-tamper (“AT”) protection is employed so that it is very difficult to reverse engineer or alter the function of electronic hardware (e.g., computer processors, integrated circuits, multi-chip modules, etc). For some commercial applications, designers often spend vast sums of money to develop a “next generation” circuit. These companies often wish to deter, or at least hamper a competitor's reverse engineering efforts. The motivation in this case is to protect valuable intellectual property. Military and Government users also have a strong interest in AT protection. When new military hardware is fielded, often the consequences of capture are not fully understood by the designer of the hardware. Similarly, the combat loss of any one of a thousand pieces of sensitive, high-tech military hardware could do irreparable damage to national security.\nMost AT is categorized as either passive or active. In each case, the intent is to delay, prevent or stop tampering and potential reverse engineering of an electronic circuit. Passive AT is currently the most widespread method of deterring an opponent from reverse engineering or spoofing an electronic circuit. Current passive AT arrangements include encapsulation or various types of conformal coatings such as epoxies. Methods to defeat common encapsulants are well documented.\nAT standards have been defined according to the Federal Information Protection Standard (FIPS) 140-2. The standard describes the requirements for four levels of protection. For the standards for multi-chip, embedded modules, Level 1 calls for standard passivation techniques (i.e., a sealing coat applied over the chip circuitry to protect it against environmental or other physical damage). The standard describes that Level 2 can be achieved using anti-tamper coatings or passive AT. Level 3 may use passive AT if tampering will likely destroy the module. Level 4 requires the use of active AT technologies.\nLayered anti-tamper arrangements are also employed in which alternating layers of passive AT with active AT yields a synergy in probing difficulty. With active AT methods, a protected circuit will take some action when unauthorized activities are detected. Any number of events can trigger a programmed circuit response. Examples of active triggering arrangements include: voltage, photon detection, acceleration, strain, thermal, chemical attack, and proximity or tamper-respondent enclosures. A tamper-respondent package can theoretically detect probing by proximity detection or by an external activity mutilating an active circuit, exterior to what is being protected. The response of an active AT circuit upon triggering is also widely variable. For example, zeroization may be employed in which critical memory cells or an entire die can be erased. Similarly, a response can trigger overwriting of some or all of a memory die. Another detection response is to physically obliterate a critical circuit element or elements."}
{"text": "For almost as long as there have been computer networks, there have been schemes which allow computers to access each other's file systems over the network in much the same manner as they access their own local file system. The first widely used remote file access protocol was Sun Microsystems' network file system (NFS), which became very popular with the rise of Unix in the mid 1980's (see B. Nowicki, “NFS: Network File System Protocol Specification,” Network Working Group RFC1094, March 1989). At about the same time, the SMB network file sharing protocol was developed by IBM for use with their PC's. Subsequent versions of SMB have become widely used on networked PC's running Microsoft Windows, and on their fileservers.\nKeeping data in networked file systems allows users to access the same data environment from different workstations on the network, and greatly simplifies system administration and the sharing of public data. For these and other reasons, it is expected that network data repositories will become widely popular among PC users as soon as typical PC network connections become fast enough to make substantial remote storage of data practical. Indeed, some Web-based services which make specific types of user data accessible from any Web browser are already popular—for example, email services and appointment calendars. Servers for individuals' Web pages also follow the network-data model.\nMany companies are offering additional Web-based services which store their data remotely, seeking new applications that will become popular. Some of these companies also offer substantial amounts of free network-based file storage. The greatest obstacle to the acceptance of these new network-based services has been slow network connections. Most computer users currently connect to the network through a telephone modem, which provides them with a connection that is about 1000 times slower than the I/O bandwidth to their local hard disk. This makes it relatively inconvenient to use remote network-based storage for most of the applications that these users now run on their local file system.\nSome companies currently sell network-based backup services to PC users. For a fee, these companies provide a combination of PC software and networked storage space that allows users to keep a copy of their most important data remotely. For privacy, the PC software encrypts user data before sending it to be stored, using the user's individual public key. Some of these companies also offer Web-based access to backed-up data. Thus far, these companies have not achieved an appreciable penetration into the PC user market. Slow network connections, the cost and effort involved in obtaining and using such services, and a low perceived benefit attached to maintaining backups of file data, have been major obstacles. For the moment, most of the Gigabytes of programs and data that users accumulate remain exclusively on their local hard disks.\nUse of network storage is also encouraged by techniques which speed up network file transfers. One such technique involves the concept of a “digital fingerprint” of a file, also called a “hash function”, a “content signature” or a “message digest” (see R. L. Rivest, “MD4 Message Digest Algorithm,” Network Working Group RFC1186, October 1990). A fingerprint is a fixed-length value obtained by mixing all of the bits of the file together in some prescribed deterministic manner—the same data always produces the same fingerprint. The fingerprint is used as a compact representative of the whole file: if two file fingerprints don't match, then the files are different. For a well designed fingerprint, the chance that any two actual files will ever have the same fingerprint can be made arbitrarily small. Such a fingerprint serves as a unique name for the file data.\nFingerprints have been used for many years to avoid unnecessary file transfers. One application of this sort has been in Bulletin Board Systems (BBSs), which have used fingerprints since the early 1990's to avoid the communication cost of uploading file data that is already present in the BBS, but associated with a different file name. Fingerprints have also been used in BBSs to conserve storage space by not storing duplicate data (for an example of both uses, see Frederick W. Kantor's Content Signature software, FWKCS, which has been in use by bulletin boards such as Channel 1 since at least 1993). These BBSs maintain a table of fingerprints for all files already present. When a new file is uploaded for storage on the BBS, its fingerprint is taken. If the BBS already contains a file with the same fingerprint (regardless of the file's name) then the duplicate data is not stored. Similarly, a client computer wishing to store data into the BBS can compute the fingerprint of the file that it wishes to send, and send that first. If a file containing this data is already present in the BBS, then the client is informed and need not send anything.\nD. A. Farber and R. D. Lachman, in U.S. Pat. No. 5,978,791 (Data processing system using substantially unique identifiers to identify data items, whereby identical data items have the same identifiers, filed October 1997) carry the idea of file fingerprints a step further, using them as the primary identifier for all data-items stored in a file system. In their scheme, not only are fingerprints used to avoid unnecessary transmission and duplicate-storage of file data (as in the BBS scheme mentioned above), but they also use fingerprints directly to gain read access to data. In this scheme, access to “licensed” data is controlled by associating explicit lists of licensees with specific data-items. Such a control mechanism doesn't scale well when applied to intellectual property protection in general. Any data-item added to the system which is copyrighted, for example, would have to have attached to it an explicit list of all users who are legally allowed to read it. Otherwise someone can give out access to the data-item to everyone that uses the file system by anonymously publishing the fingerprint of the data-item. Constructing an explicit legal-access list for each data-item is in general cumbersome, difficult and intrusive.\nFurthermore, existing schemes which use fingerprints to identify redundant data and avoid unnecessary transmission and storage depend upon the storage system being able to examine previously stored data. If users independently encrypt their data for privacy, they can't take advantage of each others data to save on transmission or on storage. If data is unencrypted, then the storage system maintainers have complete access to all user data. They may be tempted or coerced into looking at this data, and in some situations may be legally obliged to provide parts of it to third parties."}
{"text": "Audiovisual content (e.g., television programming) may often be used to transmit information to viewers. This information may include incentives, coupons, and so on. For example, a clothing store may advertise on a certain television channel and display the address, telephone number and the like of the clothing store. However, some users may have difficulty remembering portions or all of the information.\nInformation such as an address, telephone number, store name, and so on may foe stored in a quick response (QR) code. A QR code is a matrix barcode that may be readable by mobile phones with a camera, smart phones, computing devices, specialized scanners, and so on. The matrix barcode may consist of black blocks or modules arranged in a pattern against a white background. The information encoded within the matrix barcode may be text uniform resource indicator (URI), alphanumeric, numeric and other data. Matrix barcodes storing addresses and URIs may appear in magazines, on signs, buses, business cards, other objects where users may desire information, and so on. Users with, a camera phone or other mobile device equipped with the correct reader application can convert a photographic image of the matrix barcode to display text, contact information, connect to a wireless network, open a webpage in the phone's browser, and so on.\nA number of display devices, such as televisions, may be configured to present three dimensional (3-D) images and/or 3-D video. These display devices may display content that appears as a 3-D image when viewed through 3-D-specific glasses. For example, the display device may display two images superimposed on each other, where each image is projected through a different polarizing filter. The viewer may view the image through a pair of glasses having a pair of lenses with different polarized orientations. The first orientation may correspond to the first filtered image and the second orientation may correspond to the second filtered image. Thus, as the image is displayed on the screen, each eye of the viewer may perceive a different image. This creates a 3-D image (as viewed by the viewer) as each eye sees the same image, but at a slightly different angle.\nOther techniques such as anaglyphic, alternative-frame sequencing, and autosteroscopic may also be used to create a 3-D appearance on a display device. These techniques (either with or without specialized lenses) present one version of the image to a left eye of the viewer and a second version of the image to the right eye of the viewer. The two images may be the same image (or substantially the same) but be slightly offset from each other. As the human eye processes the two images together, the brain sees a 3-D image. These 3-D techniques may be used to display audiovisual content in 3-D. For example, many movies may be presented in 3-D to allow a user to have a more interactive movie-watching experience.\nThe information included in this Background section of the specification, including any references cited herein and any description or discussion thereof, is included for technical reference purposes only and is not to be regarded subject matter by which the scope of the invention as defined in the claims is to be bound."}
{"text": "1. Field of Invention\nThe field of the currently claimed embodiments of this invention relates to solution processes for producing electrodes for electro-optic devices and devices that use the electrodes.\n2. Discussion of Related Art\nThe most dominant transparent conductive film currently in use for electro-optic devices, such as liquid crystal displays, touch-screen displays and photovoltaic cells, for example, is Indium-Tin-Oxide (ITO). High transparency and low sheet resistant have ensured its wide application in many areas. Although it has been used for several decades, there is significant interest in developing alternatives to ITO because of the decreasing availability, and accompanying increase in price, of indium. Due to the high and increasing demand for transparent electrodes, and the limited supply of indium, there is currently an urgent need for new transparent electrodes.\nThere are several candidates for transparent electrodes based on materials that can potentially replace ITO. These include the use of carbon nanotubes (CNT), graphene, or thin metal films, for example. However, these candidates all have disadvantages that they compromise between optical transparency and electrical conductivity.\nRecently, some efforts have been directed to using silver nanowire (AgNW) networks to form transparent conductors. There are several issues that remain to be addressed for the large scale fabrication of AgNW films. First, a good electrical connection between crossed AgNWs is a key factor to achieve high conductivity. However, due to surfactant coating of PVP on the surface of AgNW, extra processes are often utilized to fuse the crossed AgNWs together. These processes include high temperature thermal annealing (>150 ° C.), applying extra pressure or vacuum filtering on anodized aluminum oxide (AAO) membrane substrates and HCl vapor treatment, etc. Second, strong adhesion between AgNW and substrates is necessary to obtain stable and robust AgNW fibrous films. Substrate surface modification has been used to improve the adhesion of AgNWs on the substrate. Burying AgNW into polymer films by applying pressure can also provide strong adhesion between AgNWs and substrates. Moreover, nail polish or in-situ polymerization has also been reported to improve the adhesion. However, these efforts have not provided transparent electrodes and methods of production that are adequate to replace ITO. Therefore, there remains a need for improved transparent electrodes, methods of producing the electrodes, and devices that use the electrodes."}
{"text": "1. Field of the Invention\nThe present invention relates in general to front constructions of motor vehicles, and more particularly to the front constructions of a type that is equipped with a heat exchanger support frame.\n2. Description of the Related Art\nHitherto, to simplify the process for mounting heat exchangers, such as radiator, condenser and the like, onto a front base of motor vehicle body flowing on a vehicle body assembling line, a sub-assembly system has been widely employed. That is, in this system, before being mounted to the front base of motor vehicle body, the heat exchangers are mounted to a heat exchanger support frame to constitute a so-called xe2x80x9cfront end modulexe2x80x9d, and then, the front end module is mounted to the front base of motor vehicle body flowing on the main assembling line.\nOne of known heat exchanger support frames is shaped rectangular and generally includes upper and lower horizontal members and right and left side stays or vertical members. In this known technique, for mounting the front end module, which has the heat exchangers previously mounted thereon, onto the front base of vehicle body, the side stays of the heat exchanger support frame are fixed to front ends of right and left side members of the front base by means of bolts and nuts respectively.\nHowever, in the above-mentioned technique, because of necessity of assuring the mechanical strength and matching the size thereof with the front ends of the side members of the front base of vehicle body, the side stays of the heat exchanger support frame are compelled to have an enlarged width. This however means narrowing of an air passing aperture defined by the heat exchanger support frame, which of course affects the performance of the heat exchangers held by the support member.\nIt is therefore an object of the present invention to provide a front construction of a motor vehicle body, which is free of the above-mentioned drawback.\nAccording to a first aspect of the present invention, there is provided a front construction of a motor vehicle body which comprises a front base of the vehicle body; a heat exchanger support frame having an air passing aperture which is defined by upper and lower horizontal members and two mutually spaced side stays which extend between the upper and lower horizontal members; a skirt portion integral with the lower horizontal member and extending along a lower portion of the lower horizontal member, the skirt portion facing forward with respect to the vehicle body keeping a given inclination relative to the heat exchanger support member to serve as an air guide for the air passing aperture; recessed seat portions integrally formed by the skirt portion; and fixing members which fix the recessed seat portions to the front base.\nAccording to a second aspect of the present invention, there is provided a heat exchanger support frame which is to be mounted on a front base of vehicle body. The heat exchanger support frame comprises upper and lower horizontal members; two mutually spaced side stays which extend between the upper and lower horizontal members thereby to define a generally rectangular air passing aperture therebetween; a skirt portion integral with the lower horizontal member and extending along a lower portion of the same, the skirt portion facing forward keeping a given inclination relative to the heat exchanger support frame to serve as an air guide for the air passing aperture; and recessed seat portions integrally formed by the skirt portion, the recessed seat portions causing the skirt portion to have a three dimensional structure having an increased mechanical strength."}
{"text": "Sleep is generally beneficial and restorative to a person. Therefore, it exerts a great influence on a person's quality of life. The human sleep/wake cycle generally conforms to a circadian rhythm that is regulated by a biological clock. Regular periods of sleep enable the body and mind to rejuvenate and rebuild. The body may perform various tasks during sleep, such as organizing long term memory, integrating new information, and renewing tissue and other body structures.\nLack of sleep and/or decreased sleep quality may have a number of causal factors including, e.g., respiratory disturbances, nerve or muscle disorders, and emotional conditions, such as depression and anxiety. Chronic long-term sleep-related disorders such as chronic insomnia, sleep-disordered breathing, and sleep movement disorders may significantly affect a patient's sleep quality and quality of life.\nSleep apnea, for example, is a fairly common breathing disorder characterized by periods of interrupted breathing experienced during sleep. Sleep apnea is typically classified based on its etiology. One type of sleep apnea, denoted as obstructive sleep apnea, occurs when the patient's airway is obstructed by the collapse of soft tissue in the rear of the throat. Central sleep apnea is caused by a derangement of the central nervous system control of respiration. The patient ceases to breathe when control signals from the brain to the respiratory muscles are absent or interrupted. Mixed apnea is a combination of the central and obstructive sleep apnea types. Regardless of the type of apnea people experiencing an apnea event stop breathing for a period of time. The cessation of breathing may occur repeatedly during sleep, sometimes hundreds of times a night and occasionally for a minute or longer.\nIn addition to apnea, other types of disordered breathing have been identified, including, for example, hypopnea (shallow breathing), dyspnea (labored breathing), hyperpnea (deep breathing), and tachypnea (rapid breathing). Combinations of the disordered respiratory events described above have also been observed. For example, Cheyne-Stokes respiration (CSR, which is sometimes referred to as periodic breathing) is associated with rhythmic increases and decreases in tidal volume caused by alternating periods of hyperpnea followed by apnea or hypopnea. The breathing interruptions of CSR may be associated with central apnea, or may be obstructive in nature. CSR is frequently observed in patients with congestive heart failure (CHF) and is associated with an increased risk of accelerated CHF progression."}
{"text": "The invention relates to a locking device, in particular for a door or the like of a vehicle, according to the features of the preamble of patent claim 1.\nDE 197 39 340.3 discloses a locking device which is used, in particular, for a door, a trunk lid or a tailgate or the like of a vehicle. Here, an actuating drive is provided, which transmits its movement, via a transmission element, to an element of the locking device, various functional positions of the locking device, such as a theft protection position, central locking position, etc., being set, depending on the movement of the actuating drive. Although DE 197 39 340.3 describes the operative connection of the actuating drive to the element of the locking device via the transmission element, no indications as to the activation of the actuating drive or transmission element and as to the detection of their positions can be gathered from this patent application.\nThe object on which the invention is based, therefore, is to provide a locking device which specifies not only the implementation of the functional positions, but also the activation of the actuating drive and the execution of activation."}
{"text": "There is a general need for a simple and accurate apparatus for rapidly distributing material in the form of small solid bodies simultaneously and uniformly into a plurality of spaced containers for further processing. For example, seeds sown in containers must be covered to prevent them from washing away, and to prevent the build-up of moss and algae on top of the containers. A number of devices have been developed to spread covering material, usually granite grit or basalt gravel, over the seed. These prior devices are relatively bulky, mechanically complex, and somewhat unreliable. For example, some of the prior devices employ revolving drum members, rotating discs, or shutter boxes, all of which are quite bulky and expensive, as well as not being satisfactorily precise. Consequently, there is a definite need for a simpler, less expensive, and more accurate manually operable spreader usable with various materials, which can be adapted for use with different types of containers and with different kinds of seed covering materials, or for distributing seeds themselves."}
{"text": "Thrust bearing assemblies generally include a pair of axially spaced races and a circular complement of rolling bearing elements received between pathways ground into the races. The bearing is so named because, while operating, it is under a compressive load that forces or thrusts the races axially toward one another. While the load may have a significant radial component, it does not generally reverse axial direction. Consequently, such a bearing needs no structure to prevent the races from axially separating once the bearing is installed and operating. However, such a bearing still must be shipped and handled prior to installation, and it is then subjected to forces that would axially separate the races, in the absence of some structure to retain them together.\nOne common application for thrust bearings is the suspension of a vehicle, especially steerable McPherson struts, where the provision of a rolling bearing gives superior steering and handling feel. A vehicle suspension is an application where unitization of the bearing during installation and handling on the line is desirable. Furthermore, a vehicle suspension is a harsh environment where a sealed bearing is almost essential. There are many different approaches to both sealing and unitization in steerable McPherson strut bearings disclosed in the published patents. One approach is to mount the bearing within the elastomer material of the the upper mount of the suspension strut itself so as to wrap and seal the bearing. An example may be seen in U.S. Pat. No. 4,274,655, assigned to the assignee of the present invention. Another possibility, when the bearing can be partially protected and shielded within the pad, is to use a non-rubbing labyrinth seal, as may be seen in U.S. Pat. No. 4,541,744, assigned to the assignee of the present invention. Such an approach finds utility where it is practical to mount the bearing within the elastomer pad of the upper mount, which is not always possible.\nAnother approach is to mold elastomer seals directly to one of the bearing races and snap engage the seal with the other race, thereby both sealing and unitizing the bearing. Two examples may be seen in U.S. Pat. Nos. 4,120,543 and 4,400,041, also assigned to the same assignee. However, there is a great cost disadvantage to directly molding elastomer to steel, as the molding process must be carefully controlled. There is also the disadvantage that the elasticity of the elastomer is then inevitably limited by the inelasticity of the steel to which it is bonded. By far the least expensive seal is a free molded seal, used in a bearing in which some structure other than the seal itself unitizes the bearing. An example may be seen in U.S. Pat. No. 4,497,523, where free molded O-rings compressed between the bearing races provide the seal, while separate snap engaged housing members enclose the races to unitize the bearing. A disadvantage of an O-ring is that it requires the compressive operating load to keep it in sealing engagement, and it disengages almost immediately upon any axial separation of the races. While that is not a problem when separate housings prevent the races from axially separating prior to installation, separate housings themselves present a cost penalty. In addition, although it is not visible in the scale of the drawings in U.S. Pat. No. 4,497,523, a conventional O-ring presents problems of its own. It cannot be allowed to just float freely between the races, some mounting groove must be added to one of the races to confine the O-ring and to prevent it from gathering as it rubs on the other race. U.S. Pat. No. 3,567,296 illustrates such an O-ring mounting groove. Besides necessitating an extra manufacturing step, the provision of an O-ring mounting groove in one race destroys the symmetry of the bearing, a symmetry which is itself very desirable from a manufacturing and installation standpoint. Such problems with conventional free molded O-rings have even led to designs in which the O-ring is, in effect, directly molded to one of the races, despite the cost penalty, see U.S. Pat. No. 4,462,608, also assigned to the assignee of the present invention."}
{"text": "Millions of people suffer from various types of respiratory disorders that cause abnormal breathing patterns. Examples of such respiratory disorders include sleep apnea, asthma and emphysema, each of which is briefly discussed below.\nSleep Apnea\nSleep apnea is a disease in which breathing ceases during sleep. It is associated with increased mortality and morbidity, including sudden death, heart failure, hypertension, stroke, and neuropsychiatric and behavioral disturbances.\nObstructive sleep apnea (OSA) is a very common type of sleep apnea, which often affects obese persons. This disorder occurs when a person's upper airway becomes partially or completely blocked, causing breathing to be interrupted intermittently during sleep. The blocked or obstructed breathing prevents deep stages of sleep, often resulting in daytime drowsiness.\nThere are a number of different reasons that the upper airway can become obstructed. For example, obstruction can result from upper airway abnormalities, such as enlarged tonsils or enlarged adenoids. Other upper airway abnormalities may result from obesity, an abnormally wide neck, or from an anatomically narrowed upper airway. Such abnormalities often result in excess tissue within the throat, which cause the upper airway to become blocked when throat muscles relax during sleep.\nThe most commonly prescribed treatment for obstructive sleep apnea is continuous positive airway pressure (CPAP). In CPAP treatment, the patient sleeps wearing a facial or nasal mask that delivers air pressure that stents the airway open. For many patients, CPAP therapy works well. However, many other patients find the masks very uncomfortable, preventing them from sleeping, and often causing them to remove the masks and thus avoid their treatment. Further, CPAP therapy has had no effect in some patients and has even resulted in increased episodes of apnea in some patients.\nOther conventional treatments for obstructive sleep apnea include use of anti-depressant drugs, such as protriptyline (often marketed as Vivactil™ or Triptil™). In addition to not always being effective, protriptyline has a number of side effects, such as dizziness, drowsiness, dryness of mouth, headaches, increased appetite, nausea, tiredness or weakness, unpleasant taste and weight gain. Less frequent side effect include diarrhea, heartburn, increased sweating, trouble in sleeping, and vomiting. There are also rare but more severe side effects, such as blurred vision, confusion or delirium, constipation (especially in the elderly), decreased sexual ability (more common with amoxapine and clomipramine), difficulty in speaking or swallowing, eye pain, fainting, fast or irregular heartbeat (pounding, racing, skipping) and hallucinations.\nOSA can also be treated surgically. However the costs are high, the risks are high, and the success rates vary greatly depending on the procedure and the experience and skills of the surgeon. Additionally, the surgery can affect speech and the ability to swallow.\nAnother type of sleep apnea is central sleep apnea (CSA), which is a different than OSA. CSA is a neurological disorder that results from problems with breathing control mechanisms, as opposed to airway blockage (as with OSA). CSA can occur, for example, due to the portion of the brain that controls breathing failing to signal the body to breathe. Persons suffering from CSA are often awoken during sleep, leading to sleep deprivation, and thus, daytime drowsiness.\nWhile CPAP therapy is sometime used with patients suffering from CSA, other types of mechanical ventilation are often the only treatment available to ensure continued breathing. Unfortunately, complications may result from prolonged mechanical ventilation.\nA further type of respiratory disorder that occurs during sleep is hypopnea, which is generally defined as decrease in airflow by at least 50% for ten seconds or more. In this document, hypopnea will be considered a type of sleep apnea, even though shallow breathing continues during hypopnea.\nDiagnosis of sleep apnea is not simple because there can be many different reasons for disturbed sleep. While there are various tests for evaluating a person for sleep apnea, most are very expensive and sometime inconclusive. Typically a patient is required to sleep overnight at a sleep treatment center, so that the patient can be monitored and diagnosed. Such overnight stays are typically very expensive and are sometimes inconclusive because the patient either did not exhibit their normal sleep patterns or because the patient does not experience sleep apnea every night. Additionally, since treatments for OSA and CSA are often different, it is important that the type of sleep apnea be accurately diagnosed.\nAs is evident from the above description, there is the need for improved systems and methods for diagnosing sleep apnea and treating sleep apnea. This will allow appropriate treatments to be prescribed. Further, even after a treatment is being provided, it is important to chronically monitor a person's breathing in order to determine if specific treatments are working, to determine if treatments can be fine tuned and/or to determine if alternative treatments should be used.\nAsthma\nAsthma is a disease of the respiratory system, which includes a person's nose, mouth, trachea, lungs, and air tubes or airways that connect the nose and mouth with the lungs (these tubes are called bronchi and bronchioles). When a person breathes, the muscles that are wrapped around the airways are very loose and relaxed. Since the lining inside the airways is very thin, the relaxed state of the muscles allows the airways to open widely so that air easily passes in and out of the small air sacs (called alveoli) that make up the lungs. Breathing includes the moving of air in and out of the lungs.\nDuring an asthma attack, the muscles around the airways tighten, or “spasm” and the lining inside the airways swell or thicken, and become blocked with thick mucous. This makes the airways much thinner than usual so it is harder to move air in and out of the air sacs, making it difficult to breathe. During an asthma attack, it is actually harder to breathe out than it is to breathe in. This means that during an asthma attack, it takes much longer to breathe out (expire) than it does to breathe in (inspire). Since it is so hard to breathe out during an asthma attack, more and more air gets trapped inside the lungs making a person feel like they can't breathe in or out.\nThere are different types of asthma, including allergy-induced asthma and activity-induced asthma. Additionally, there are different severities of asthma, including mild, intermittent, moderate and severe asthma. Different types of asthma and different severities of asthma may be treated in different manners. Accordingly, it is desirable to be able to obtain information that is useful in diagnosing the type and/or severity of a person's asthmatic condition. Further, even after a person's asthma has been diagnosed, the proper treatment needs to be provided. Treatments include medication, altering diet and/or altering activities. However, even if two different people have the same symptoms, they may not react in the same way to the same treatment. That is, not all types of treatments work well with all people. Accordingly, even after a treatment regimen has been prescribed, it is important to chronically monitor a person's breathing in order to determine if specific treatments are working, to determine if treatments can be fine tuned and/or to determine if alternative treatments should be used.\nEmphysema\nEmphysema is an abnormal condition of the lungs, which belongs to a group of conditions called Chronic Obstructive Pulmonary Disease (COPD). Emphysema, which generally develops gradually over a number of years, typically results in an increased shortness of breath, not relieved by sitting upright, and a chronic cough. The shortness of breath is generally worse on exertions such as walking up a flight of stairs or eating a big meal. Eventually, breathing becomes harder and harder causing fatigue, weakness, and weight loss. Emphysema can result for various reasons, such as chronic bronchitis, or chronic irritation to the lungs by dust, pollution, and cigarette smoking. Such chronic irritations cause the tiny air sacs in the lungs to loose their elasticity and become distended, causing inhaled air to become trapped in the air sacs. As with most respiratory disorders, there is no cure for emphysema but there are numerous treatments that can help relieve the shortness of breath.\nThere are different causes and severities for emphysema, each of which can be treated in multiple different manners. Accordingly, it is desirable to be able to obtain information that is useful in diagnosing the type and/or severity of the emphysema. Further, even after a treatment is being provided, it is important to chronically monitor a person's breathing in order to determine if specific treatments are working, to determine if treatments can be fine tuned and/or to determine if alternative treatments should be used.\nAs can be appreciated from the above description, it would be useful to provide a means for temporarily and/or chronically monitoring various breathing disorders in order to properly diagnose the disorders, treat the disorders, and monitor the treatment of the disorders."}
{"text": "Recorded audio and motion picture media has been an aspect of society since the days of Thomas Edison. At the start of the 20th century there was wide distribution of recorded audio media (cylinders and records) and motion picture media (nickelodeons and movies), but both technologies were still in their infancy. In the late 1920s motion pictures were combined with audio on a mass-market basis, followed by color motion pictures with audio. Radio broadcasting gradually evolved into a largely advertising-supported form of broadcast mass-market audio media. When a television (TV) broadcast standard was established in the mid-1940s, television joined radio as a form of broadcast mass-market media bringing previously recorded or live motion pictures into the home.\nBy the middle of the 20th century, a large percentage of US homes had phonograph record players for playing recorded audio media, a radio to receive live broadcast audio, and a television set to play live broadcast audio/video (A/V) media. Very often these 3 “media players” (record player, radio and TV) were combined into one cabinet sharing common speakers that became the “media center” for the home. Although the media choices were limited to the consumer, the media “ecosystem” was quite stable. Most consumers knew how to use the “media players” and were able to enjoy the full extent of their capabilities. At the same time, the publishers of the media (largely the motion picture and televisions studios, and the music companies) were able to distribute their media both to theaters and to the home without suffering from widespread piracy or “second sales”, i.e., the resale of used media. Typically publishers do not derive revenue from second sales, and as such, it reduces revenue that publishers might otherwise derive from the buyer of used media for new sales. Although there certainly were used records sold during the middle of the 20th century, such sales did not have a large impact on record publishers because, unlike a motion picture or video program—which is typically watched once or only a few times by an adult—a music track may be listened to hundreds or even thousands of times. So, music media is far less “perishable” (i.e., it has lasting value to an adult consumer) than motion picture/video media. Once a record was purchased, if the consumer liked the music, the consumer was likely to keep it a long time.\nFrom the middle of the 20th century through the present day, the media ecosystem has undergone a series of radical changes, both to the benefit and the detriment of consumers and publishers. With the widespread introduction of audio recorders, especially cassette tapes with high-quality stereo sound, there certainly was a higher degree of consumer convenience. But it also marked the beginning of what is now a widespread practice with consumer media: piracy. Certainly, many consumers used the cassette tapes for taping their own records purely for convenience, but increasingly consumers (e.g., students in a dormitory with ready access to each others' record collections) would make pirated copies. Also, consumers would tape music played over the radio rather than buying a record or tape from the publisher.\nThe advent of the consumer VCR led to even more consumer convenience, since now a VCR could be set to record a TV show which could be watched at a later time, and it also led to the creation of the video rental business, where movies as well as TV programming could be accessed on an “on demand” basis. The rapid development of mass-market home media devices since the mid-1980s has led to an unprecedented level of choice and convenience for the consumer, and also has led to a rapid expansion of the media publishing market.\nToday, consumers are faced with a plethora of media choices as well as a plethora of media devices, many of which are tied to particular forms of media or particular publishers. An avid consumer of media may have a stack of devices connected to TVs and computers in various rooms of the house, resulting in a “rat's nest” of cables to one or more TV sets and/or personal computers (PCs) as well as a group of remote controls. (In the context of the present application, the term “personal computer” or “PC” refers to any sort of computer suitable for us in the home or office, including a desktop, a Macintosh® or other non-Windows computers, Windows-compatible devices, Unix variations, laptops, etc.) These devices may include a video game console, VCR, DVD player, audio surround-sound processor/amplifier, satellite set-top box, cable TV set-top box, etc. And, for an avid consumer, there may be multiple similar-function devices because of compatibility issues. For example, a consumer may own both a HD-DVD and a Blu-ray DVD player, or both a Microsoft Xbox® and a Sony Playstation® video game system. Indeed, because of incompatibility of some games across versions of game consoles, the consumer may own both an XBox and a later version, such as an Xbox 360®. Frequently, consumers are befuddled as to which video input and which remote to use. Even after a disc is placed into the correct player (e.g., DVD, HD-DVD, Blu-ray, Xbox or Playstation), the video and audio input is selected for that the device, and the correct remote control is found, the consumer is still faced with technical challenges. For example, in the case of a wide-screen DVD, the user may need to first determine and then set the correct aspect ratio on his TV or monitor screen (e.g., 4:3, Full, Zoom, Wide Zoom, Cinema Wide, etc.). Similarly, the user may need to first determine and then set the correct audio surround sound system format (e.g., AC-3, Dolby Digital, DTS, etc.). Often times, the consumer is unaware that they may not be enjoying the media content to the full capability of their television or audio system (e.g., watching a movie squashed at the wrong aspect ratio, or listening to audio in stereo rather than in surround sound).\nIncreasingly, Internet-based media devices have been added to the stack of devices. Audio devices like the Sonos® Digital Music system stream audio directly from the Internet. Likewise, devices like the Slingbox™ entertainment player record video and stream it through a home network or out through the Internet where it can be watched remotely on a PC. And Internet Protocol Television (IPTV) services offer cable TV-like services through Digital Subscriber Line (DSL) or other home Internet connections. There have also been recent efforts to integrate multiple media functions into a single device, such as the Moxi® Media Center and PCs running Windows XP Media Center Edition. While each of these devices offers an element of convenience for the functions that it performs, each lacks ubiquitous and simple access to most media. Further, such devices frequently cost hundreds of dollars to manufacture, often because of the need for expensive processing and/or local storage. Additionally, these modern consumer electronic devices typically consume a great deal of power, even while idle, which means they are expensive over time and wasteful of energy resources. For example, a device may continue to operate if the consumer neglects to turn it off or switches to a different video input. And, because none of the devices is a complete solution, it must be integrated with the other stack of devices in the home, which still leaves the user with a rat's nest of wires and a sea of remote controls.\nFurthermore, when many newer Internet-based devices do work properly, they typically offer media in a more generic form than it might otherwise be available. For example, devices that stream video through the Internet often stream just the video material, not the interactive “extras” that often accompany DVDs, like the “making of” videos, games, or director's commentary. This is due to the fact that frequently the interactive material is produced in a particular format intended for a particular device that handles interactivity locally. For example, each of DVD, HD-DVDs and Blu-ray discs have their own particular interactive format. Any home media device or local computer that might be developed to support all of the popular formats would require a level of sophistication and flexibility that would likely make it prohibitively expensive and complex for the consumer to operate.\nAdding to the problem, if a new format were introduced later in the future the local device may not have the hardware capability to support the new format, which would mean that the consumer would have to purchase an upgraded local media device. For example, if higher-resolution video or stereoscopic video (e.g., one video stream for each eye) were introduced at a later date, the local device may not have the computational capability to decode the video, or it may not have the hardware to output the video in the new format (e.g., assuming stereoscopy is achieved through 120 fps video synchronized with shuttered glasses, with 60 fps delivered to each eye, if the consumer's video hardware can only support 60 fps video, this option would be unavailable absent an upgraded hardware purchase).\nThe issue of media device obsolescence and complexity is a serious problem when it comes to sophisticated interactive media, especially video games.\nModern video game applications are largely divided into four major non-portable hardware platforms: Sony PlayStation® 1, 2 and 3 (PS1, PS2, and PS3); Microsoft Xbox® and Xbox 360®; and Nintendo Gamecube® and Wii™; and PC-based games. Each of these platforms is different than the others so that games written to run on one platform usually do not run on another platform. There may also be compatibility problems from one generation of device to the next. Even though the majority of software game developers create software games that are designed independent of a particular platform, in order to run a particular game on a specific platform a proprietary layer of software (frequently called a “game development engine”) is needed to adapt the game for use on a specific platform. Each platform is sold to the consumer as a “console” (i.e., a standalone box attached to a TV or monitor/speakers) or it is a PC itself. Typically, the video games are sold on optical media such as a Blu-ray DVD, DVD-ROM or CD-ROM, which contains the video game embodied as a sophisticated real-time software application. As home broadband speeds have increased, video games are becoming increasingly available for download.\nThe specificity requirements to achieve platform-compatibility with video game software is extremely exacting due to the real-time nature and high computational requirements of advanced video games. For example, one might expect full game compatibility from one generation to the next of video games (e.g., from XBox to XBox 360, or from Playstation 2 (“PS2”) to Playstation 3 (“PS3”), just as there is general compatibility of productivity applications (e.g., Microsoft Word) from one PC to another with a faster processing unit or core. However, this is not the case with video games. Because the video game manufacturers typically are seeking the highest possible performance for a given price point when a video game generation is released, dramatic architectural changes to the system are frequently made such that many games written for the prior generation system do not work on the later generation system. For example, XBox was based upon the x86-family of processors, whereas XBox 360 was based upon a PowerPC-family.\nTechniques can be utilized to emulate a prior architecture, but given that video games are real-time applications, it is often unfeasible to achieve the exact same behavior in an emulation. This is a detriment to the consumer, the video game console manufacturer and the video game software publisher. For the consumer, it means the necessity of keeping both an old and new generation of video game consoles hooked up to the TV to be able to play all games. For the console manufacturer it means cost associated with emulation and slower adoption of new consoles. And for the publisher it means that multiple versions of new games may have to be released in order to reach all potential consumers—not only releasing a version for each brand of video game (e.g., XBox, Playstation), but often a version for each version of a given brand (e.g., PS2 and PS3). For example, a separate version of Electronic Arts' “Madden NFL 08” was developed for XBox, XBox 360, PS2, PS3, Gamecube, Wii, and PC, among other platforms.\nPortable devices, such as cellular (“cell”) phones and portable media players also present challenges to game developers. Increasingly such devices are connected to wireless data networks and are able to download video games. But, there are a wide variety of cell phones and media devices in the market, with a wide range of different display resolutions and computing capabilities. Also, because such devices typically have power consumption, cost and weight constraints, they typically lack advanced graphics acceleration hardware like a Graphics Processing Unit (“GPU”), such as devices made by NVIDIA of Santa Clara, Calif. Consequently, game software developers typically develop a given game title simultaneously for many different types of portable devices. A user may find that a given game title is not available for his particular cell phone or portable media player.\nIn the case of home game consoles, hardware platform manufacturers typically charge a royalty to the software game developers for the ability to publish a game on their platform. Cell phone wireless carriers also typically charge a royalty to the game publisher to download a game into the cell phone. In the case of PC games, there is no royalty paid to publish games, but game developers typically face high costs due to the higher customer service burden to support the wide range of PC configurations and installation issues that may arise. Also, PCs typically present less barriers to the piracy of game software since they are readily reprogrammable by a technically-knowledgeable user and games can be more easily pirated and more easily distributed (e.g., through the Internet). Thus, for a software game developer, there are costs and disadvantages in publishing on game consoles, cell phones and PCs.\nFor game publishers of console and PC software, costs do not end there. To distribute games through retail channels, publishers charge a wholesale price below the selling price for the retailer to have a profit margin. The publisher also typically has to pay the cost of manufacturing and distributing the physical media holding the game. The publisher is also frequently charged a “price protection fee” by the retailer to cover possible contingencies such as where the game does not sell, or if the game's price is reduced, or if the retailer must refund part or all of the wholesale price and/or take the game back from a buyer. Additionally, retailers also typically charge fees to publishers to help market the games in advertising flyers. Furthermore, retailers are increasingly buying back games from users who have finished playing them, and then sell them as used games, typically sharing none of the used game revenue with the game publisher. Adding to the cost burden placed upon game publishers is the fact that games are often pirated and distributed through the Internet for users to download and make free copies.\nAs Internet broadband speeds have been increasing and broadband connectivity has become more widespread in the US and worldwide, particularly to the home and to Internet “cafes” where Internet-connected PCs are rented, games are increasingly being distributed via downloads to PCs or consoles. Also, broadband connections are increasingly used for playing multiplayer and massively multiplayer online games (both of which are referred to in the present disclosure by the acronym “MMOG”). These changes mitigate some of the costs and issues associated with retail distribution. Downloading online games addresses some of the disadvantages to game publishers in that distribution costs typically are less and there are little or no costs from unsold media. But downloaded games are still subject to piracy, and because of their size (often many gigabytes in size) they can take a very long time to download. In addition, multiple games can fill up small disk drives, such as those sold with portable computers or with video game consoles. However, to the extent games or MMOGs require an online connection for the game to be playable, the piracy problem is mitigated since the user is usually required to have a valid user account. Unlike linear media (e.g., video and music) which can be copied by a camera shooting video of the display screen or a microphone recording audio from the speakers, each video game experience is unique, and can not be copied using simple video/audio recording. Thus, even in regions where copyright laws are not strongly enforced and piracy is rampant, MMOGs can be shielded from piracy and therefore a business can be supported. For example, Vivendi SA's “World of Warcraft” MMOG has been successfully deployed without suffering from piracy throughout the world. And many online or MMOG games, such as Linden Lab's “Second Life” MMOG generate revenue for the games' operators through economic models built into the games where assets can be bought, sold, and even created using online tools. Thus, mechanisms in addition to conventional game software purchases or subscriptions can be used to pay for the use of online games.\nWhile piracy can be often mitigated due to the nature of online or MMOGs, online game operator still face remaining challenges. Many games require substantial local (i.e., in-home) processing resources for online or MMOGs to work properly. If a user has a low performance local computer (e.g., one without a GPU, such as a low-end laptop), he may not be able to play the game. Additionally, as game consoles age, they fall further behind the state-of-the-art and may not be able to handle more advanced games. Even assuming the user's local PC is able to handle the computational requirements of a game, there are often installation complexities. There may be driver incompatibilities (e.g., if a new game is downloaded, it may install a new version of a graphics driver that renders a previously-installed game, reliant upon an old version of the graphics driver, inoperable). A console may run out of local disk space as more games are downloaded. Complex games typically receive downloaded patches over time from the game developer as bugs are found and fixed, or if modifications are made to the game (e.g., if the game developer finds that a level of the game is too hard or too easy to play). These patches require new downloads. But sometimes not all users complete downloading of all the patches. Other times, the downloaded patches introduce other compatibility or disk space consumption issues.\nAlso, during game play, large data downloads may be required to provide graphics or behavioral information to the local PC or console. For example, if the user enters a room in a MMOG and encounters a scene or a character made up of graphics data or with behaviors that are not available on the user's local machine, then that scene or character's data must be downloaded. This may result in a substantial delay during game play if the Internet connection is not fast enough. And, if the encountered scene or character requires storage space or computational capability beyond that of the local PC or console, it can create a situation where the user can not proceed in the game, or must continue with reduced-quality graphics. Thus, online or MMOG games often limit their storage and/or computational complexity requirements. Additionally, they often limit the amount of data transfers during the game. Online or MMOG games may also narrow the market of users that can play the games.\nFurthermore, technically-knowledgeable users are increasingly reverse-engineering local copies of games and modifying the games so that they can cheat. The cheats maybe as simple as making a button press repeat faster than is humanly possible (e.g., so as to shoot a gun very rapidly). In games that support in-game asset transactions the cheating can reach a level of sophistication that results in fraudulent transactions involving assets of actual economic value. When an online or MMOGs economic model is based on such asset transactions, this can result in substantial detrimental consequences to the game operators.\nThe cost of developing a new game has grown as PCs and consoles are able to produce increasingly sophisticated games (e.g., with more realistic graphics, such as real-time ray-tracing, and more realistic behaviors, such as real-time physics simulation). In the early days of the video game industry, video game development was a very similar process to application software development; that is, most of the development cost was in the development of the software, as opposed to the development of the graphical, audio, and behavioral elements or “assets”, such as those that may be developed for a motion picture with extensive special effects. Today, many sophisticated video game development efforts more closely resemble special effects-rich motion picture development than software development. For instance, many video games provide simulations of 3-D worlds, and generate increasingly photorealistic (i.e., computer graphics that seem as realistic as live action imagery shot photographically) characters, props, and environments. One of the most challenging aspects of photorealistic game development is creating a computer-generated human face that is indistinguishable from a live action human face. Facial capture technologies such Contour™ Reality Capture developed by Mova of San Francisco, Calif. captures and tracks the precise geometry of a performer's face at high resolution while it is in motion. This technology allows a 3D face to be rendered on a PC or game console that is virtually indistinguishable from a captured live action face. Capturing and rendering a “photoreal” human face precisely is useful in several respects. First, highly recognizable celebrities or athletes are often used in video games (often hired at a high cost), and imperfections may be apparent to the user, making the viewing experience distracting or unpleasant. Frequently, a high degree of detail is required to achieve a high degree of photorealism—requiring the rendering of a large number of polygons and high-resolution textures, potentially with the polygons and/or textures changing on a frame-by-frame basis as the face moves.\nWhen high polygon-count scenes with detailed textures change rapidly, the PC or game console supporting the game may not have sufficient RAM to store enough polygon and texture data for the required number of animation frames generated in the game segment. Further, the single optical drive or single disk drive typically available on a PC or game console is usually much slower than the RAM, and typically can not keep up with the maximum data rate that the GPU can accept in rendering polygons and textures. Current games typically load most of the polygons and textures into RAM, which means that a given scene is largely limited in complexity and duration by the capacity of the RAM. In the case of facial animation, for example, this may limit a PC or a game console to either a low resolution face that is not photoreal, or to a photoreal face that can only be animated for a limited number of frames, before the game pauses, and loads polygons and textures (and other data) for more frames.\nWatching a progress bar move slowly across the screen as a PC or console displays a message similar to “Loading . . . ” is accepted as an inherent drawback by today's users of complex video games. The delay while the next scene loads from the disk (“disk” herein, unless otherwise qualified, refers to non-volatile optical or magnetic media, as well non-disk media such as semiconductor “Flash” memory) can take several seconds or even several minutes. This is a waste of time and can be quite frustrating to a game player. As previously discussed, much or all of the delay may be due to the load time for polygon, textures or other data from a disk, but it also may be the case that part of the load time is spent while the processor and/or GPU in the PC or console prepares data for the scene. For example, a soccer video game may allow the players to choose among a large number of players, teams, stadiums and weather conditions. So, depending on what particular combination is chosen, different polygons, textures and other data (collectively “objects”) may be required for the scene (e.g., different teams have different colors and patterns on their uniforms). It may be possible to enumerate many or all of the various permutations and pre-compute many or all of the objects in advance and store the objects on the disk used to store the game. But, if the number of permutations is large, the amount of storage required for all of the objects may be too large to fit on the disk (or too impractical to download). Thus, existing PC and console systems are typically constrained in both the complexity and play duration of given scenes and suffer from long load times for complex scenes.\nAnother significant limitation with prior art video game systems and application software systems is that they are increasingly using large databases, e.g., of 3D objects such as polygons and textures, that need to be loaded into the PC or game console for processing. As discussed above, such databases can take a long time to load when stored locally on a disk. Load time, however, is usually far more severe if the database is stored a remote location and is accessed through the Internet. In such a situation it may take minutes, hours, or even days to download a large database. Further, such databases are often created a great expense (e.g., a 3D model of a detailed tall-masted sailing ship for use in a game, movie, or historical documentary) and are intended for sale to the local end-user. However, the database is at risk of being pirated once it has been downloaded to the local user. In many cases, a user wants to download a database simply for the sake of evaluating it to see if it suits the user's needs (e.g., if a 3D costume for a game character has a satisfactory appearance or look when the user performs a particular move). A long load time can be a deterrent for the user evaluating the 3D database before deciding to make a purchase.\nSimilar issues occur in MMOGs, particularly as games that allow users to utilize increasingly customized characters. For a PC or game console to display a character it needs to have access to the database of 3D geometry (polygons, textures, etc.) as well as behaviors (e.g., if the character has a shield, whether the shield is strong enough to deflect a spear or not) for that character. Typically, when a MMOG is first played by a user, a large number of databases for characters are already available with the initial copy of the game, which is available locally on the game's optical disk or downloaded to a disk. But, as the game progresses, if the user encounters a character or object whose database is not available locally (e.g., if another user has created a customized character), before that character or object can be displayed, its database must be downloaded. This can result in a substantial delay of the game.\nGiven the sophistication and complexity of video games, another challenge for video game developers and publishers with prior art video game consoles, is that it frequently takes 2 to 3 years to develop a video game at a cost of tens of millions of dollars. Given that new video game console platforms are introduced at a rate of roughly once every five years, game developers need to start development work on those games years in advance of the release of the new game console in order to have video games available concurrently when the new platform is released. Several consoles from competing manufactures are sometimes released around the same time (e.g., within a year or two of each other), but what remains to be seen is the popularity of each console, e.g., which console will produce the largest video game software sales. For example, in a recent console cycle, the Microsoft XBox 360, the Sony Playstation 3, and the Nintendo Wii were scheduled to be introduced around the same general timeframe. But years before the introductions the game developers essentially had to “place their bets” on which console platforms would be more successful than others, and devote their development resources accordingly. Motion picture production companies also have to apportion their limited production resources based on what they estimate to be the likely success of a movie well in advance of the release of the movie. Given the growing level of investment required for video games, game production is increasingly becoming like motion picture production, and game production companies routinely devote their production resources based on their estimate of the future success of a particular video game. But, unlike they motion picture companies, this bet is not simply based on the success of the production itself; rather, it is predicated on the success of the game console the game is intended to run on. Releasing the game on multiple consoles at once may mitigate the risk, but this additional effort increases cost, and frequently delays the actual release of the game.\nApplication software and user environments on PCs are becoming more computationally intensive, dynamic and interactive, not only to make them more visually appealing to users, but also to make them more useful and intuitive. For example, both the new Windows Vista™ operating system and successive versions of the Macintosh® operating system incorporate visual animation effects. Advanced graphics tools such as Maya™ from Autodesk, Inc., provide very sophisticated 3D rendering and animation capability which push the limits of state-of-the-art CPUs and GPUs. However, the computational requirements of these new tools create a number of practical issues for users and software developers of such products.\nSince the visual display of an operating system (OS) must work on a wide range of classes of computers—including prior-generation computers no longer sold, but still upgradeable with the new OS—the OS graphical requirements are limited to a large degree by a least common denominator of computers that the OS is targeted for, which typically includes computers that do not include a GPU. This severely limits the graphics capability of the OS. Furthermore, battery-powered portable computers (e.g., laptops) limit the visual display capability since high computational activity in a CPU or GPU typically results in higher power consumption and shorter battery life. Portable computers typically include software that automatically lowers processor activity to reduce power consumption when the processor is not utilized. In some computer models the user may lower processor activity manually. For example, Sony's VGN-SZ280P laptop contains a switch labeled “Stamina” on one side (for low performance, more battery life) and “Speed” on the other (for high performance, less battery life). An OS running on a portable computer must be able to function usably even in the event the computer is running at a fraction of its peak performance capability. Thus, OS graphics performance often remains far below the state-of-the-art available computational capability.\nHigh-end computationally-intense applications like Maya are frequently sold with the expectation that they will be used on high-performance PCs. This typically establishes a much higher performance, and more expensive and less portable, least common denominator requirement. As a consequence, such applications have a much more limited target audience than a general purpose OS (or general purpose productivity application, like Microsoft Office) and typically sell in much lower volume than general purpose OS software or general purpose application software. The potential audience is further limited because often times it is difficult for a prospective user to try out such computationally-intense applications in advance. For example, suppose a student wants to learn how to use Maya or a potential buyer already knowledgeable about such applications wants to try out Maya before making the investment in the purchase (which may well involve also buying a high-end computer capable of running Maya). While either the student or the potential buyer could download, or get a physical media copy of, a demo version of Maya, if they lack a computer capable of running Maya to its full potential (e.g., handling a complex 3D scene), then they will be unable to make an fully-informed assessment of the product. This substantially limits the audience for such high-end applications. It also contributes to a high selling price since the development cost is usually amortized across a much smaller number of purchases than those of a general-purpose application.\nHigh-priced applications also create more incentive for individuals and businesses to use pirated copies of the application software. As a result, high-end application software suffers from rampant piracy, despite significant efforts by publishers of such software to mitigate such piracy through various techniques. Still, even when using pirated high-end applications, users cannot obviate the need to invest in expensive state-of-the-art PCs to run the pirated copies. So, while they may obtain use of a software application for a fraction of its actual retail price, users of pirated software are still required to purchase or obtain an expensive PC in order to fully utilize the application.\nThe same is true for users of high-performance pirated video games. Although pirates may get the games at fraction of their actual price, they are still required to purchase expensive computing hardware (e.g., a GPU-enhanced PC, or a high-end video game console like the XBox 360) needed to properly play the game. Given that video games are typically a pastime for consumers, the additional cost for a high-end video game system can be prohibitive. This situation is worse in countries (e.g., China) where the average annual income of workers currently is quite low relative to that of the United States. As a result, a much smaller percentage of the population owns a high-end video game system or a high-end PC. In such countries, “Internet cafes”, in which users pay a fee to use a computer connected to the Internet, are quite common. Frequently, such Internet cafes have older model or low-end PCs without high performance features, such as a GPU, which might otherwise enable players to play computationally-intensive video games. This is a key factor in the success of games that run on low-end PCs, such as Vivendi's “World of Warcraft” which is highly successful in China, and is commonly played in Internet cafes there. In contrast, a computationally-intensive game, like “Second Life” is much less likely to be playable on a PC installed in a Chinese Internet café. Such games are virtually inaccessible to users who only have access to low-performance PCs in Internet cafes.\nBarriers also exist for users who are considering purchasing a video game and would first like to try out a demonstration version of the game by downloading the demo through the Internet to their home. A video game demo is often a full-fledged version of the game with some features disabled, or with limits placed on the amount of game play. This may involve a long process (perhaps hours) of downloading gigabytes of data before the game can be installed and executed on either a PC or a console. In the case of a PC, it may also involve figuring out what special drivers are needed (e.g., DirectX or OpenGL drivers) for the game, downloading the correct version, installing them, and then determining whether the PC is capable of playing the game. This latter step may involve determining whether the PC has enough processing (CPU and GPU) capability, sufficient RAM, and a compatible OS (e.g., some games run on Windows XP, but not Vista). Thus, after a long process of attempting to run a video game demo, the user may well find out that the video game demo can't be possibly played, given the user's PC configuration. Worse, once the user has downloaded new drivers in order to try the demo, these driver versions may be incompatible with other games or applications the user uses regularly on the PC, thus the installation of a demo may render previously operable games or applications inoperable. Not only are these barriers frustrating for the user, but they create barriers for video game software publishers and video game developers to market their games.\nAnother problem that results in economic inefficiency has to do with the fact that given PC or game console is usually designed to accommodate a certain level of performance requirement for applications and/or games. For example, some PCs have more or less RAM, slower or faster CPUs, and slower or faster GPUs, if they have a GPUs at all. Some games or applications make take advantage of the full computing power of a given PC or console, while many games or applications do not. If a user's choice of game or application falls short of the peak performance capabilities of the local PC or console, then the user may have wasted money on the PC or console for unutilized features. In the case of a console, the console manufacturer may have paid more than was necessary to subsidize the console cost.\nAnother problem that exists in the marketing and enjoyment of video games involves allowing a user to watch others playing games before the user commits to the purchase of that game. Several prior art approaches exist for the recording of video games for replay at a later time. For example, U.S. Pat. No. 5,558,339 teaches recording game state information, including game controller actions, during “gameplay” in the video game client computer (owned by the same or different user). This state information can be used at a later time to replay some or all of the game action on a video game client computer (e.g., PC or console). A significant drawback to this approach is that for a user to view the recorded game, the user must possess a video game client computer capable of playing the game and must have the video game application running on that computer, such that the gameplay is identical when the recorded game state is replayed. Beyond that, the video game application has to be written in such a way that there is no possible execution difference between the recorded game and the played back game.\nFor example, game graphics are generally computed on a frame-by-frame basis. For many games, the game logic sometimes may take shorter or longer than one frame time to compute the graphics displayed for the next frame, depending on whether the scene is particularly complex, or if there are other delays that slow down execution (e.g., on a PC, another process may be running that takes away CPU cycles from the game applications). In such a game, a “threshold” frame that is computed in slightly less than one frame time (say a few CPU clock cycles less) can eventually occur. When that same scene is computed again using the exact same game state information, it could easily take a few CPU clock cycles more than one frame time (e.g., if an internal CPU bus is slightly out of phase with the an external DRAM bus and it introduces a few CPU cycle times of delay, even if there is no large delay from another process taking away milliseconds of CPU time from game processing). Therefore, when the game is played back the frame gets calculated in two frame times rather than a single frame time. Some behaviors are based on how often the game calculates a new frame (e.g., when the game samples the input from the game controllers). While the game is played, this discrepancy in the time reference for different behaviors does not impact game play, but it can result in the played-back game producing a different result. For example, if a basketball's ballistics are calculated at a steady 60 fps rate, but the game controller input is sampled based on rate of computed frames, the rate of computed frames may be 53 fps when the game was recorded, but 52 fps when the game is replayed, which can make the difference between whether the basketball is blocked from going into the basket or not, resulting in a different outcome. Thus, using game state to record video games requires very careful game software design to ensure that the replay, using the same game state information, produces the exact same outcome.\nAnother prior art approach for recording video game is to simply record the video output of a PC or video game system (e.g., to a VCR, DVD recorder, or to a video capture board on a PC). The video then can be rewound and replayed, or alternatively, the recorded video uploaded to the Internet, typically after being compressed. A disadvantage to this approach is that when a 3D game sequence is played back, the user is limited to viewing the sequence from only the point of view from which the sequence was recorded. In other words, the user cannot change the point of view of the scene.\nFurther, when compressed video of a recorded game sequence played on a home PC or game console is made available to other users through the Internet, even if the video is compressed in real-time, it may be impossible to upload the compressed video in real-time to the Internet. The reason why is because many homes in the world that are connected to the Internet have highly asymmetric broadband connections (e.g., DSL and cable modems typically have far higher downstream bandwidth than upstream bandwidth). Compressed high resolution video sequences often have higher bandwidths than the upstream bandwidth capacity of the network, making them impossible to upload in real-time. Thus, there would be a significant delay after the game sequence is played (perhaps minutes or even hours) before another user on the Internet would be able to view the game. Although this delay is tolerable in certain situations (e.g., to watch a game player's accomplishments that occurred at a prior time), it eliminates the ability to watch a game live (e.g., a basketball tournament, played by champion players) or with “instant replay” capability as the game is played live.\nAnother prior art approach allows a viewer with a television receiver to watch video games live, but only under the control of the television production crew. Some television channels, in both the US and in other countries provide video game viewing channels, where the television viewing audience is able to watch certain video game users (e.g., top-rated players playing in tournaments) on video game channels. This is accomplished by having the video output of the video game systems (PCs and/or consoles) fed into the video distribution and processing equipment for the television channel. This is not unlike when the television channel is broadcasting a live basketball game in which several cameras provide live feeds from different angles around the basketball court. The television channel then is able to make use of their video/audio processing and effects equipment to manipulate the output from the various video game systems. For example, the television channel can overlay text on top of the video from a video game that indicates the status of different players (just as they might overlay text during a live basketball game), and the television channel can overdub audio from a commentator who can discuss the action occurring during the games. Additionally, the video game output can be combined with cameras recording video of the actual players of the games (e.g., showing their emotional response to the game).\nOne problem with this approach is that such live video feeds must be available to the television channel's video distribution and processing equipment in real-time in order for it to have the excitement of a live broadcast. As previously discussed, however, this is often impossible when the video game system is running from the home, especially if part of the broadcast includes live video from a camera that is capturing real-world video of the game player. Further, in a tournament situation, there is a concern that an in-home gamer may modify the game and cheat, as previously described. For these reasons, such video game broadcasts on television channels are often arranged with players and video game systems aggregated at a common location (e.g., at a television studio or in an arena) where the television production equipment can accept video feeds from multiple video game systems and potentially live cameras.\nAlthough such prior art video game television channels can provide a very exciting presentation to the television viewing audience that is an experience akin to a live sporting event, e.g., with the video game players presented as “athletes”, both in terms of their actions in the video game world, and in terms of their actions in the real world, these video game systems are often limited to situations where players are in close physical proximity to one another. And, since television channels are broadcasted, each broadcasted channel can only show one video stream, which is selected by the television channel's production crew. Because of these limitations and the high cost of broadcast time, production equipment and production crews, such television channels typically only show top-rated players playing in top tournaments.\nAdditionally, a given television channel broadcasting a full-screen image of a video game to the entire television viewing audience shows only one video game at a time. This severely limits a television viewer's choices. For example, a television viewer may not be interested in the game(s) shown at a given time. Another viewer may only be interested in watching the game play of a particular player that is not featured by the television channel at a given time. In other cases, a viewer may only be interested in watching a how an expert player handles a particular level in a game. Still other viewers may wish to control the viewpoint that a video game is seen from, which is different from that chosen by the production team, etc. In short, a television viewer may have a myriad of preferences in watching video games that are not accommodated by the particular broadcast of a television network, even if several different television channels are available. For all of the aforementioned reasons, prior art video game television channels have significant limitations in presenting video games to television viewers.\nAnother drawback of prior art video games systems and application software systems is that they are complex, and commonly suffer from errors, crashes and/or unintended and undesired behaviors (collectively, “bugs”). Although games and applications typically go through a debugging and tuning process (frequently called “Software Quality Assurance” or SQA) before release, almost invariably once the game or application is released to a wide audience in the field bugs crop up. Unfortunately, it is difficult for the software developer to identify and track down many of the bugs after release. It can be difficult for software developers to become aware of bugs. Even when they learn about a bug, there may only be a limited amount of information available to them to identify what caused the bug. For example, a user may call up a game developer's customer service line and leave a message stating that when playing the game, the screen started to flash, then changed to a solid blue color and the PC froze. That provides the SQA team with very little information useful in tracking down a bug. Some games or applications that are connected online can sometimes provide more information in certain cases. For example, a “watchdog” process can sometimes be used to monitor the game or application for “crashes”. The watchdog process can gather statistics about the status of the game or applications process (e.g., the status of the stack, of the memory usage, how far the game or applications has progressed, etc.) when it crashes and then upload that information to the SQA team via the Internet. But in a complex game or application, such information can take a very long time to decipher in order to accurately determine what the user was doing at the time of the crash. Even then, it may be impossible to determine what sequence of events led to the crash.\nYet another problem associated with PCs and game consoles is that they are subject to service issues which greatly inconvenience the consumer. Service issues also impact the manufacturer of the PC or game console since they typically are required to send a special box to safely ship the broken PC or console, and then incur the cost of repair if the PC or console is in warranty. The game or application software publisher can also be impacted by the loss of sales (or online service use) by PCs and/or consoles being in a state of repair.\nFIG. 1 illustrates a prior art video gaming system such as a Sony Playstation® 3, Microsoft Xbox 360®, Nintendo Wii™, Windows-based personal computer or Apple Macintosh. Each of these systems includes a central processing unit (CPU) for executing program code, typically a graphical processing unit (GPU) for performing advanced graphical operations, and multiple forms of input/output (I/O) for communicating with external devices and users. For simplicity, these components are shown combined together as a single unit 100. The prior art video gaming system of FIG. 1 also is shown including an optical media drive 104 (e.g., a DVD-ROM drive); a hard drive 103 for storing video game program code and data; a network connection 105 for playing multi-player games, for downloading games, patches, demos or other media; a random access memory (RAM) 101 for storing program code currently being executed by the CPU/GPU 100; a game controller 106 for receiving input commands from the user during gameplay; and a display device 102 (e.g., a SDTV/HDTV or a computer monitor).\nThe prior art system shown in FIG. 1 suffers from several limitations. First, optical drives 104 and hard drives 103 tend to have much slower access speeds as compared to that of RAM 101. When working directly through RAM 101, the CPU/GPU 100 can, in practice, process far more polygons per second than is possible when the program code and data is read directly off of hard drive 103 or optical drive 104 due to the fact that RAM 101 generally has much higher bandwidth and does not suffer from the relatively long seek delays of disc mechanisms. But only a limited amount of RAM is provided in these prior art systems (e.g., 256-512 Mbytes). Therefore, a “Loading . . . ” sequence in which RAM 101 is periodically filled up with the data for the next scene of the video game is often required.\nSome systems attempt to overlap the loading of the program code concurrently with the gameplay, but this can only be done when there is a known sequence of events (e.g., if a car is driving down a road, the geometry for the approaching buildings on the roadside can be loaded while the car is driving). For complex and/or rapid scene changes, this type of overlapping usually does not work. For example, in the case where the user is in the midst of a battle and RAM 101 is completely filled with data representing the objects within view at that moment, if the user moves the view rapidly to the left to view objects that are not presently loaded in RAM 101, a discontinuity in the action will result since there not be enough time to load the new objects from Hard Drive 103 or Optical Media 104 into RAM 101.\nAnother problem with the system of FIG. 1 arises due to limitations in the storage capacity of hard drives 103 and optical media 104. Although disk storage devices can be manufactured with a relatively large storage capacity (e.g., 50 gigabytes or more), they still do not provide enough storage capacity for certain scenarios encountered in current video games. For example, as previously mentioned, a soccer video game might allow the user to choose among dozens of teams, players and stadiums throughout the world. For each team, each player and each stadium a large number of texture maps and environment maps are needed to characterize the 3D surfaces in the world (e.g., each team has a unique jersey, with each requiring a unique texture map).\nOne technique used to address this latter problem is for the game to pre-compute texture and environment maps once they are selected by the user. This may involve a number of computationally-intensive processes, including decompressing images, 3D mapping, shading, organizing data structures, etc. As a result, there may be a delay for the user while the video game is performing these calculations. On way to reduce this delay, in principle, is to perform all of these computations—including every permutation of team, player roster, and stadium—when the game was originally developed. The released version of the game would then include all of this pre-processed data stored on optical media 104, or on one or more servers on the Internet with just the selected pre-processed data for a given team, player roster, stadium selection downloaded through the Internet to hard drive 103 when the user makes a selection. As a practical matter, however, such pre-loaded data of every permutation possible in game play could easily be terabytes of data, which is far in excess of the capacity of today's optical media devices. Furthermore, the data for a given team, player roster, stadium selection could easily be hundreds of megabytes of data or more. With a home network connection of, say, 10 Mbps, it would take longer to download this data through network connection 105 than it would to compute the data locally.\nThus, the prior art game architecture shown in FIG. 1 subjects the user to significant delays between major scene transitions of complex games.\nAnother problem with prior art approaches such as that shown in FIG. 1 is that over the years video games tend to become more advanced and require more CPU/GPU processing power. Thus, even assuming an unlimited amount of RAM, video games hardware requirements go beyond the peak level of processing power available in these systems. As a result, users are required to upgrade gaming hardware every few years to keep pace (or play newer games at lower quality levels). One consequence of the trend to ever more advanced video games is that video game playing machines for home use are typically economically inefficient because their cost is usually determined by the requirements of the highest performance game they can support. For example, an XBox 360 might be used to play a game like “Gears of War”, which demands a high performance CPU, GPU, and hundreds of megabytes of RAM, or the XBox 360 might be used to play Pac Man, a game from the 1970s that requires only kilobytes of RAM and a very low performance CPU. Indeed, an XBox 360 has enough computing power to host many simultaneous Pac Man games at once.\nVideo games machines are typically turned off for most of the hours of a week. According to a July 2006 Nielsen Entertainment study of active gamers 13 years and older, on average, active gamers spend fourteen hours/week playing console video games, or just 12% of the total hours in a week. This means that the average video game console is idle 88% of the time, which is an inefficient use of an expensive resource. This is particularly significant given that video game consoles are often subsidized by the manufacturer to bring down the purchase price (with the expectation that the subsidy will be earned back by royalties from future video game software purchases).\nVideo game consoles also incur costs associated with almost any consumer electronic device. For instance, the electronics and mechanisms of the systems need to be housed in an enclosure. The manufacturer needs to offer a service warranty. The retailer who sells the system needs to collect a margin on either the sale of the system and/or on the sale of video game software. All of these factors add to the cost of the video game console, which must either be subsidized by the manufacturer, passed along to the consumer, or both.\nIn addition, piracy is a major problem for the video game industry. The security mechanisms utilized on virtually every major video gaming system have been “cracked” over the years, resulting in unauthorized copying of video games. For example, the Xbox 360 security system was cracked in July 2006 and users are now able to download illegal copies online. Games that are downloadable (e.g., games for the PC or the Mac) are particularly vulnerable to piracy. In certain regions of the world where piracy is weakly policed there is essentially no viable market for standalone video game software because users can buy pirated copies as readily as legal copies for a tiny fraction of the cost. Also, in many parts of the world the cost of a game console is such a high percentage of income that even if piracy were controlled, few people could afford a state-of-the-art gaming system.\nIn addition, the used game market reduces revenue for the video game industry. When a user has become tired of a game, they can sell the game to a store which will resell the game to other users. This unauthorized but common practice significantly reduces revenues of game publishers. Similarly, a reduction in sales on the order of 50% commonly occurs when there is a platform transition every few years. This is because users stop buying games for the older platforms when they know that the newer version platform is about to be released (e.g., when Playstation 3 is about to be released, users stop buying Playstation 2 games). Combined, the loss of sales and increased development costs associated with the new platforms can have a very significant adverse impact on the profitability of game developers.\nNew game consoles are also very expensive. The Xbox 360, the Nintendo Wii, and the Sony Playstation 3 all retail for hundreds of dollars. High powered personal computer gaming systems can cost up to $8000. This represents a significant investment for users, particularly considering that the hardware becomes obsolete after a few years and the fact that many systems are purchased for children.\nOne approach to the foregoing problems is online gaming in which the gaming program code and data are hosted on a server and delivered to client machines on-demand as compressed video and audio streamed over a digital broadband network. Some companies such as G-Cluster in Finland (now a subsidiary of Japan's SOFTBANK Broadmedia) currently provide these services online. Similar gaming services have become available in local networks, such as those within hotels and offered by DSL and cable television providers. A major drawback of these systems is the problem of latency, i.e., the time it takes for a signal to travel to and from the game server, which is typically located in an operator's “head-end”. Fast action video games (also known as “twitch” video games) require very low latency between the time the user performs an action with the game controller and the time the display screen is updated showing the result of the user action. Low latency is needed so that the user has the perception that the game is responding “instantly”. Users may be satisfied with different latency intervals depending on the type of game and the skill level of the user. For example, 100 ms of latency may be tolerable for a slow casual game (like backgammon) or a slow-action role playing game, but in a fast action game a latency in excess of 70 or 80 ms may cause the user to perform more poorly in the game, and thus is unacceptable. For instance, in a game that requires fast reaction time there is a sharp decline in accuracy as latency increases from 50 to 100 ms.\nWhen a game or application server is installed in a nearby, controlled network environment, or one where the network path to the user is predictable and/or can tolerate bandwidth peaks, it is far easier to control latency, both in terms of maximum latency and in terms of the consistency of the latency (e.g., so the user observes steady motion from digital video streaming through the network). Such level of control can be achieved between a cable TV network head-end to a cable TV subscriber's home, or from a DSL central office to DSL subscriber's home, or in a commercial office Local Area Network (LAN) environment from a server or a user. Also, it is possible to obtain specially-graded point-to-point private connections between businesses which have guaranteed bandwidth and latency. But in a game or application system that hosts games in a server center connected to the general Internet and then streams compressed video to the user through a broadband connection, latency is incurred from many factors, resulting in severe limitations in the deployment of prior art systems.\nIn a typical broadband-connected home, a user may have a DSL or cable modem for broadband service. Such broadband services commonly incur as much as a 25 ms round-trip latency (and at times more) between the user's home and the general Internet. In addition, there are round-trip latencies incurred from routing data through the Internet to a server center. The latency through the Internet varies based on the route that the data is given and the delays it incurs as it is routed. In addition to routing delays, round-trip latency is also incurred due to the speed of light traveling through the optical fiber that interconnects most of the Internet. For example, for each 1000 miles, approximately 22 ms is incurred in round-trip latency due to the speed of light through the optical fiber and other overhead.\nAdditional latency can occur due to the data rate of the data streamed through the Internet. For example, if a user has DSL service that is sold as “6 Mbps DSL service”, in practice, the user will probably get less than 5 Mbps of downstream throughput at best, and will likely see the connection degrade periodically due to various factors such as congestion during peak load times at the Digital Subscriber Line Access Multiplexer (DSLAM). A similar issue can occur reducing a the data rate of a cable modem is used for a connection sold as “6 Mbps cable modem service” to far less than that, if there is congestion in the local shared coaxial cable looped through the neighborhood, or elsewhere in the cable modem system network. If data packets at a steady rate of 4 Mbps are streamed as one-way in User Datagram Protocol (UDP) format from a server center through such connections, if everything is working well, the data packets will pass through without incurring additional latency, but if there is congestion (or other impediments) and only 3.5 Mbps is available to stream data to the user, then in a typical situation either packets will be dropped, resulting in lost data, or packets will queue up at the point of congestion, until they can be sent, thereby introducing additional latency. Different points of congestion have different queuing capacity to hold delayed packets, so in some cases packets that can't make it through the congestion are dropped immediately. In other cases, several megabits of data are queued up and eventually be sent. But, in almost all cases, queues at points of congestion have capacity limits, and once those limits are exceeded, the queues will overflow and packets will be dropped. Thus, to avoid incurring additional latency (or worse, loss of packets), it is necessary to avoid exceeding the data rate capacity from the game or application server to the user.\nLatency is also incurred by the time required to compress video in the server and decompress video in the client device. Latency is further incurred while a video game running on a server is calculating the next frame to be displayed. Currently available video compression algorithms suffer from either high data rates or high latency. For example, motion JPEG is an intraframe-only lossy compression algorithm that is characterized by low-latency. Each frame of video is compressed independently of each other frame of video. When a client device receives a frame of compressed motion JPEG video, it can immediately decompress the frame and display it, resulting in very low latency. But because each frame is compressed separately, the algorithm is unable to exploit similarities between successive frames, and as a result intraframe-only video compression algorithms suffer from very high data rates. For example, 60 fps (frames per second) 640×480 motion JPEG video may require 40 Mbps (megabits per second) or more of data. Such high data rates for such low resolution video windows would be prohibitively expensive in many broadband applications (and certainly for most consumer Internet-based applications). Further, because each frame is compressed independently, artifacts in the frames that may result from the lossy compression are likely to appear in different places in successive frames. This can results in what appears to the viewer as a moving visual artifacts when the video is decompressed.\nOther compression algorithms, such as MPEG2, H.264 or VC9 from Microsoft Corporation as they are used in prior art configurations, can achieve high compression ratios, but at the cost of high latency. Such algorithms utilize interframe as well as intraframe compression. Periodically, such algorithms perform an intraframe-only compression of a frame. Such a frame is known as a key frame (typically referred to as an “I” frame). Then, these algorithms typically compare the I frame with both prior frames and successive frames. Rather than compressing the prior frames and successive frames independently, the algorithm determines what has changed in the image from the I frame to the prior and successive frames, and then stores those changes as what are called “B” frames for the changes preceding the I frame and “P” frames for the changes following the I frame. This results in much lower data rates than intraframe-only compression. But, it typically comes at the cost of higher latency. An I frame is typically much larger than a B or P frame (often 10 times larger), and as a result, it takes proportionately longer to transmit at a given data rate.\nConsider, for example, a situation where the I frames are 10× the size of B and P frames, and there are 29 B frames+30 P frames=59 interframes for every single I intraframe, or 60 frames total for each “Group of Frames” (GOP). So, at 60 fps, there is 1 60-frame GOP each second. Suppose the transmission channel has a maximum data rate of 2 Mbps. To achieve the highest quality video in the channel, the compression algorithm would produce a 2 Mbps data stream, and given the above ratios, this would result in 2 Megabits (Mb)/(59+10)=30,394 bits per intraframe and 303,935 bits per I frame. When the compressed video stream is received by the decompression algorithm, in order for the video to play steadily, each frame needs to decompressed and displayed at a regular interval (e.g., 60 fps). To achieve this result, if any frame is subject to transmission latency, all of the frames need to be delayed by at least that latency, so the worst-case frame latency will define the latency for every video frame. The I frames introduce the longest transmission latencies since they are largest, and an entire I frame would have to be received before the I frame could be decompressed and displayed (or any interframe dependent on the I frame). Given that the channel data rate is 2 Mbps, it will take 303,935/2 Mb=145 ms to transmit an I frame.\nAn interframe video compression system as described above using a large percentage of the bandwidth of the transmission channel will be subject to long latencies due to the large size of an I frame relative to the average size of a frame. Or, to put it another way, while prior art interframe compression algorithms achieve a lower average per-frame data rate than intraframe-only compression algorithms (e.g., 2 Mbps vs. 40 Mbps), they still suffer from a high peak per-frame data rate (e.g., 303,935*60=18.2 Mbps) because of the large I frames. Bear in mind, though that the above analysis assumes that the P and B frames are all much smaller than the I frames. While this is generally true, it is not true for frames with high image complexity uncorrelated with the prior frame, high motion, or scene changes. In such situations, the P or B frames can become as large as I frames (if a P or B frame gets larger than an I frame, a sophisticated compression algorithm will typically “force” an I frame and replace the P or B frame with an I frame). So, I frame-sized data rate peaks can occur at any moment in a digital video stream. Thus, with compressed video, when the average video data rate approaches data rate capacity of the transmission channels (as is frequently the case, given the high data rate demands for video) the high peak data rates from I frames or large P or B frames result in a high frame latency.\nOf course, the above discussion only characterizes the compression algorithm latency created by large B, P or I frames in a GOP. If B frames are used, the latency will be even higher. The reason why is because before a B frame can be displayed, all of the B frames after the B frame and the I frame must be received. Thus, in a group of picture (GOP) sequence such as BBBBBIPPPPPBBBBBIPPPPP, where there are 5 B frames before each I frame, the first B frame can not be displayed by the video decompressor until the subsequent B frames and I frame are received. So, if video is being streamed at 60 fps (i.e., 16.67 ms/frame), before the first B frame can be decompressed, five B frames and the I frame will take 16.67*6=100 ms to receive, no matter how fast the channel bandwidth is, and this is with just 5 B frames. Compressed video sequences with 30 B frames are quite common. And, at a low channel bandwidth like 2 Mbps, the latency impact caused by the size of the I frame is largely additive to the latency impact due to waiting for B frames to arrive. Thus, on a 2 Mbps channel, with a large number of B frames it is quite easy to exceed 500 ms of latency or more using prior art video compression technology. If B frames are not used (at the cost of a lower compression ratio for given quality level), the B frame latency is not incurred, but the latency caused by the peak frame sizes, described above, is still incurred.\nThe problem is exacerbated by very the nature of many video games. Video compression algorithms utilizing the GOP structure described above have been largely optimized for use with live video or motion picture material intended for passive viewing. Typically, the camera (whether a real camera, or a virtual camera in the case of a computer-generated animation) and scene is relatively steady, simply because if the camera or scene moves around too jerkily, the video or movie material is (a) typically unpleasant to watch and (b) if it is being watched, usually the viewer is not closely following the action when the camera jerks around suddenly (e.g., if the camera is bumped when shooting a child blowing out the candles on a birthday cake and suddenly jerks away from the cake and back again, the viewers are typically focused on the child and the cake, and disregard the brief interruption when the camera suddenly moves). In the case of a video interview, or a video teleconference, the camera may be held in a fixed position and not move at all, resulting in very few data peaks at all. But 3D high action video games are characterized by constant motion (e.g., consider a 3D racing, where the entire frame is in rapid motion for the duration of the race, or consider first-person shooters, where the virtual camera is constantly moving around jerkily). Such video games can result in frame sequences with large and frequent peaks where the user may need to clearly see what is happening during those sudden motions. As such, compression artifacts are far less tolerable in 3D high action video games. Thus, the video output of many video games, by their nature, produces a compressed video stream with very high and frequent peaks.\nGiven that users of fast-action video games have little tolerance for high latency, and given all of the above causes of latency, to date there have been limitations to server-hosted video games that stream video on the Internet. Further, users of applications that require a high degree of interactivity suffer from similar limitations if the applications are hosted on the general Internet and stream video. Such services require a network configuration in which the hosting servers are set up directly in a head end (in the case of cable broadband) or the central office (in the case of Digital Subscriber Lines (DSL)), or within a LAN (or on a specially-graded private connection) in a commercial setting, so that the route and distance from the client device to the server is controlled to minimize latency and peaks can be accommodated without incurring latency. LANs (typically rated at 100 Mbps-1 Gbps) and leased lines with adequate bandwidth typically can support peak bandwidth requirements (e.g., 18 Mbps peak bandwidth is a small fraction of a 100 Mbps LAN capacity).\nPeak bandwidth requirements can also be accommodated by residential broadband infrastructure if special accommodations are made. For example, on a cable TV system, digital video traffic can be given dedicated bandwidth which can handle peaks, such as large I frames. And, on a DSL system, a higher speed DSL modem can be provisioned, allowing for high peaks, or a specially-graded connection can provisioned which can handle a higher data rates. But, conventional cable modem and DSL infrastructure attached to the general Internet have far less tolerance for peak bandwidth requirements for compressed video. So, online services that host video games or applications in server centers a long distance from the client devices, and then stream the compressed video output over the Internet through conventional residential broadband connections suffer from significant latency and peak bandwidth limitations—particularly with respect to games and applications which require very low latency (e.g., first person shooters and other multi-user, interactive action games, or applications requiring a fast response time)."}
{"text": "1. Field of the Invention\nThe present invention relates to specific polymeric dispersants and stable pigment dispersions and inkjet inks made with the specific polymeric dispersant.\n2. Description of the Related Art\nA dispersant is a substance for promoting the formation and stabilization of a dispersion of pigment particles in a dispersion medium. Dispersants are generally surface-active materials having an anionic, cationic or non-ionic structure. The presence of a dispersant substantially reduces the dispersing energy required. Dispersed pigment particles may have a tendency to re-agglomerate after the dispersing operation, due to mutual attraction forces. The use of dispersants counteracts this re-agglomeration tendency of the pigment particles.\nThe dispersant has to meet particularly high requirements when used for inkjet inks. Inadequate dispersing manifests itself as increased viscosity in liquid systems, loss of brilliance and/or hue shifts. Moreover, particularly good dispersion of the pigment particles is required to ensure unimpeded passage of pigment particles through the nozzles of a print head, which are usually only a few micrometers in diameter. In addition, pigment particle agglomeration and the associated blockage of the printer nozzles has to be avoided in the standby periods of the printer.\nPolymeric dispersants typically contain in one part of the molecule so-called anchor groups, which adsorb onto the pigments to be dispersed. In a spatially separate part of the molecule, polymeric dispersants have polymer chains compatible with the dispersion medium, thus stabilizing the pigment particles in the dispersion medium. Typical polymeric dispersants include graft copolymer and block copolymer dispersants.\nIn aqueous inkjet inks, the polymeric dispersants generally contain hydrophobic anchor groups exhibiting a high affinity for the pigment surface and hydrophilic polymer chains for stabilizing the pigments in the aqueous dispersion medium.\nThe preparation of good thermally stable dispersions with submicron particles is more difficult for non-aqueous inkjet inks, such as solvent based, oil based and radiation curable inkjet inks. The pigments are especially difficult to disperse when they have a non-polar surface.\nThese problems have led to the design of very specific polymeric dispersants wherein the anchor groups are pigment derivatives. For example, WO 2007/006635 (AGFA GRAPHICS) discloses pigment dispersions including a colour pigment and a polymeric dispersant having via a linking group covalently linked to its polymeric backbone at least one pending chromophore group which has a molecular weight smaller than 90% of the molecular weight of the colour pigment.\nAnother approach for dispersing pigments with non-polar surfaces in non-aqueous dispersion media is changing the surface to a more polar surface by addition of compounds known as dispersion synergists. A dispersion synergist is a compound that promotes the adsorption of the polymeric dispersant on the surface of the pigment. It is often suggested that the synergist should possess the pigment structure substituted by one or more sulphonic acid groups, carboxylic acid groups or ammonium salts thereof. Examples of these dispersion synergists are given in, for example, WO 2007/060254 (AGFA GRAPHICS), EP 1790697 A (AGFA GRAPHICS) and EP 1790698 A (AGFA GRAPHICS).\nHowever, both approaches lead to a considerable higher cost of the pigment dispersion. It is desirable to make stable pigment dispersions not requiring dispersion synergists or polymeric dispersants wherein the anchor groups are pigment derivatives. For consistent image quality, inkjet inks require a dispersion stability capable of dealing with high temperatures (above 60° C.) during transport of the ink to a customer, jetting at elevated temperatures and changes in the dispersion medium of the inkjet ink during use, for example, evaporation of solvent and increasing concentrations of humectants, penetrants and other additives.\nU.S. Pat. No. 4,119,615 (TEXACO) discloses thermoplastic adhesive compositions comprising a polyoxamide formed by reacting a polyoxypropylene polyamine selected from the group consisting of diamines, triamines and mixtures thereof, and having an average molecular weight of about 190 to about 3,000 with oxalic acid to form a liquid prepolymer, said prepolymer being further reacted with a diamine. U.S. Pat. No. 4,119,615 (TEXACO) is silent on pigment dispersions.\nU.S. Pat. No. 5,286,267 (TEXACO) discloses a polyether hydroxyethylaminoethyl oxalamide as gasoline detergent additive. U.S. Pat. No. 5,286,267 (TEXACO) is also silent on pigment dispersions.\nDE 3608153 (BASF) discloses a coating composition for wood containing medium- to long-oil alkyd resins as binder and low molecular weight oxalic acid diamide derivatives as an UV-absorber.\nTherefore, it is highly desirable to manufacture low cost non-aqueous pigment dispersions, especially pigmented inkjet inks, exhibiting a high dispersion quality and stability."}
{"text": "1. Field of the Invention\nThe invention pertains to electro-optic devices for switching unpolarized or polarized optical signals and more particularly to two-state optical switches which couple preselected pairs of terminals in each switch state.\n2. Description of the Prior Art\nLarge information transmission capacity, immunity to electromagnetic interference, and freedom from ground loop problems ideally suit optical transmission systems for linking distributed computers and computer controlled industrial system components. These optical transmission systems utilize optical fibers to serially link a multiplicity of optical repeater stations. Power failure at one of the serially linked stations, however, may interrupt the data chain and cause the entire system to fail. To prevent such a catastrophe, a fail safe optical switch is employed at each repeater and component station which operates to bypass that station when local power is lost. These fail safe switches must possess low insertion loss properties, provide high isolation between the input and the output optical fibers during the \"Power On\" mode, and must function effectively with unpolarized light coupled thereto from a multi-mode fiber.\nOne switch design to accomplish this bypass function employs mirrors positioned on a metallic bar that is held in place between input and output optical fibers by an electromagnet. Light entering the input fiber is coupled to the station receiver by reflection from a mirror to an optical fiber coupled to the receiver+s detector. The optical signal from the station transmitter is coupled to an optical fiber wherefrom it illuminates a second mirror and deflected therefrom to an output fiber to propagate therealong to the next station. When the power fails, the electromagnet is deenergized and a spring, attached to the metallic bar, pulls the metal bar-mirror combination from the path between the input and output optical fibers, thus allowing light to be directly coupled therebetween. This system provides adequate switching performance, but does not possess the reliability inherent in an electro-optic switch. Electro-optic switches of the prior art, though possessing sufficiently low insertion losses, do not exhibit sufficiently high isolation between the input and output optical fibers during the energized switch mode. One such switch is described in the technical paper \"Electrically Controlled Optical Switch for Multi-Mode Fiber Applications\" by R. E. Wagner and J. Cheng, published in Applied Optics Sept. 1, 1980, Volume 19, Number 17. While this switch possesses sufficiently low insertion loss properties, it requires a plurality of costly polarizing beam splitters and extremely tight component tolerances to achieve crosstalk ratios in the energized state of -20 dB at best, a crosstalk level that is not sufficiently low to meet system requirements."}
{"text": "1) Field of the Invention\nThe present invention relates to a process for preparing a p-type thermoelectric material that is used as a material for thermoelectric devices utilizing Peltier effect or Seebeck effect.\n2) Description of the Related Art\nA thermoelectric device utilizing Peltier effect or Seebeck effect can be used in a variety of utilities such as elements for heating and cooling and for temperature control, thermoelectric electricity generators and the like.\nThe thermoelectric device is made from a thermoelectric material. Performance properties of the thermoelectric material are evaluated by the thermoelectric figure of merit Z that is obtained through the following equation,Z=α2/(ρ·κ)in which Z[1/K] is figure of merit, α[μV/K] is Seebeck coefficient, κ [mW/cm·K] is thermal conductivity, and ρ[mΩ·cm] is specific resistance.\nIf the figure of merit Z is higher, it means that the thermoelectric material has higher performance. To increase the figure of merit Z of the thermoelectric material, the Seebeck coefficient α may be increased and/or the specific resistance ρ and the thermal conductivity κ may be decreased.\nMeanwhile, p-type and n-type thermoelectric materials obtained by adding a suitable dopant to an alloy containing at least two elements selected from the group consisting of bismuth (Bi), tellurium (Te), selenium (Se), antimony (Sb), and sulfur (S) are known to have higher figure of merit.\nOne of the methods for preparing this thermoelectric material is to mix and melt a powder obtained by weighing prescribed amounts of powders of Bi, Te, Se, Sb, S, and a prescribed amount of dopant, to obtain an alloy ingot, to grind the alloy ingot to obtain an alloy powder, and sinter the powder.\nThe sintering may be hot press sintering, normal pressure sintering, vacuum sintering, gas pressure sintering, plasma sintering, hot isostatic pressing (HIP) and the like. Inter alla, hot press sintering which hardly produces crack due to cleavage and affords a sintered material excellent in the mechanical strength is effective.\nThe figure of merit Z of a p-type thermoelectric material prepared by hot press sintering is around 3.0×10−3K−1. Generally, it is preferable that the figure of merit Z is high. For example, for use in an electric cooler or use in cooling a central processing unit (CPU) of a personal computer, a p-type thermoelectric material with the performance of 3.0×10−3K−1 may be used without practical problem.\nHowever, a p-type thermoelectric material with the figure of merit Z of 3.0×10−3K−1 is not sufficient in many cases. These cases include an optical communication semiconductor laser requiring temperature control at a precision of 0.1° C. or lower in order to control an oscillation wavelength, use in an electric constant temperature bath necessary for temperature control of solution such as a photoresist solution, a plating solution and various surface treating solutions in a semiconductor manufacturing step, constant temperature control and temperature test for various materials and parts, and temperature control of a culture solution in gene or microorganism culturing, and use in precise temperature control in an ultraprecise air temperature and humidity controlling apparatus at a semiconductor manufacturing step. Thus, a p-type thermoelectric material with higher figure of merit Z is in demand.\nHowever, it is difficult to obtain a p-type thermoelectric material having a higher figure of merit Z with the conventional process. Therefore, a process for manufacturing a p-type thermoelectric material with higher figure of merit Z is in demand."}
{"text": "The present invention relates to a monitoring status display device, a monitoring status display method, and a monitoring status display program.\nA system operation management needs to not only constantly monitor resources and the operating statuses of services, which are monitoring subjects, but also, in case of problems, promptly identify the problems and restore the system. The following Non-Patent Literature 1 discloses an integrated management tool for consolidating an entire system by collecting pieces of monitoring information from the monitoring subjects into a single integrated console. Non-Patent Literature 2 described below discloses a system operation monitoring tool that presents a list of monitored items defined based on SLA (Service Level Agreement).\nNon-Patent Literature 1: Integrated System Operation Management JP1 Version 9, Catalog by Hitachi, Ltd., November 2009, P. 5-6, Monitoring, “Searched on Feb. 10, 2010.”\nURL:http://www.hitachi.co.jp/Prod/comp/soft1/downlad/catalog/ca/ca740.pdf\nNon-Patent Literature 2: “Service Level Management using Tivoli Service Level Advisor,” Apr. 7, 2005, IBM Japan, P. 11, Customer Status Report (overall), “Searched on Feb. 10, 2010.”\nURL:http://www-06.ibm.com/jp/services/itil/pdf/events/2005040704.pdf\nThe integrated management tool described in Non-Patent Literature 1 displays a group of monitoring subjects in the form of a tree structure. The system operation monitoring tool described in Non-Patent Literature 2 displays a table showing, for example, the monitored items, the times/dates of monitoring, and the presence/absence of problems.\nTherefore, the more the monitoring subjects, the larger the tree and the table. This causes an administrator to repeatedly scroll or change the screen until the administrator can identify a monitoring subject having a problem. It requires a long time to identify a problem. The administrator might miss the problems while repeating such operation."}
{"text": "The present inventors have recognized that it is desirable to have the ability to modify video signals so they carry additional information (data) which can be received and used by a receiving device (e.g., a TV set or display device but which, for security reasons, is not recordable by standard video recorders. It is believed that this problem of video data conveyance without recordability has not been earlier recognized. For instance, the VBI (vertical blanking interval) of video signals is often used to carry (non-video) data such as closed captioning, but a standard VHS type VCR does record the VBI and also this data. However, current digital video recorders do not record the VBI, and it is believed that future digital recorders may record only, for instance, part of the amplitude of the VBI."}
{"text": "The present invention relates to a printing apparatus and method suitable for printing a digitized image at an altered resolution without loss of image quality.\nImprovements in image-processing technology have enabled facsimile machines, personal computers, workstations, and similar devices to exchange and process a wide variety of images. For example, an image received by a facsimile machine or television receiver may be captured into the memory of a personal computer, processed in the computer, then printed on a printer connected to the computer.\nA problem that arises is that different devices generate digitized images at different resolutions, and printers also have different resolutions. Thus a personal computer may receive an image having a resolution of six hundred dots per inch (600 DPI), and have to print the image on an electrophotographic printer with a resolution of 300 DPI. The simple solution of converting the resolution by skipping alternate dots and lines leads to serious image degradation.\nJapanese Patent Kokai Publication No. 71940/1994 discloses a printer that solves this problem by doubling the vertical resolution while halving the horizontal resolution. Each raster line in the original 600-DPI image yields two 300-DPI printed lines, with dots combined by AND logic in one line and exclusive-OR logic in the other line. This solution works well when the resolution ratio is 2:1, but does not apply to other resolution ratios.\nNon-integer resolution ratios are particularly troublesome. When a personal computer's facsimile application software receives image data having a resolution of 200 DPI, for example, and prints the image on a 300-DPI printer, the resolution must be increased by a factor of 1.5. Conventionally, this is done by repeating every second dot and every second line, but that leads to uneven printing of fine detail, such as the fine lines that form characters in text. Moreover, if the image contains bit-mapped data generated from a natural image by dithering, the repeating of every second line and dot can generate moire patterns that greatly impair the appearance of the image.\nAnother known technique accordingly divides an image into a text part and a non-text part, converts the resolution of the text part by a mapping method that takes line thicknesses into account, and converts the resolution of the non-text part by a method such as error diffusion that avoids moire patterns. Text and non-text parts must be accurately identified, however, because if the wrong method is applied, the appearance is the printed output is considerably worsened. Complex and expensive processing circuits are required to recognize text and non-text parts of images with reasonable accuracy, and even these circuits become unreliable when text is superimposed on a natural image.\nThe problems described above occur not only with electrophotographic printers, but also with thermal printers and ink-jet printers which, like electrophotographic printers, form images from dots."}
{"text": "This invention relates to an imagable clean release laminate construction, and in particular to a construction in which a releasable adhesive secures together two or more sheets which can be imaged through the use of chemical carbonless techniques.\nNumerous types of businesses produce a wide variety of labels, nameplates, identification cards tags, forms, and the like in large quantities. Each such product typically has specific indicia printed thereon by high speed printing devices. It is desirable to produce such products on a continuous web, with individual labels, forms, or the like being adhesively attached to the web.\nIn the past, pressure sensitive adhesives have been utilized to secure the products to a continuous web during printing. Typically, labels or the like were attached to a silicone release liner by a pressure sensitive adhesive. Customarily when the pressure sensitive adhesive label was removed from the release liner, the label was attached to another surface by means of the tacky adhesive backing on the label.\nHowever, in some instances, it is desirable to be able to have both surfaces of a label or the like, and also the surface of the web from which the label is released, to be clean and non-tacky. Such products are known as clean release products and are typically used as sew on labels, nameplates, and identification cards.\nPreviously, such clean release products have been made through the use of so-called releasable or fugitive adhesives. These releasable adhesives exhibit cohesive (internal) failure instead of the usual adhesive/adherend failure exhibited by typical pressure sensitive adhesives. These releasable adhesives are compositions which form relatively weak bonds but which adhere well to paper and other surfaces. Failure occurs within the releasable adhesive film itself when a layer of a laminate is subjected to peeling forces. This class of adhesives can be formulated to provide a range of cohesive strengths.\nWhen cohesive failure occurs in such releasable adhesives a thin film of adhesive remains on the back surface of the face sheet, on the face of the backing sheet or both. However the film does not exhibit tack and has no adverse effect on the appearance or handling characteristics of the product. Moreover the surface can be readily printed on.\nIn some instances it is desirable to have a product which can be imprinted and which has the added function of providing one or more duplicate copies of the printed information. For example, Singer, U.S. Pat. No. 3,383,121, teaches a self-adhesive copy label which is affixed to a backer sheet by means of a pressure sensitive adhesive. Singer utilizes a multilayer structure which includes a self-contained carbonless sheet which produces the duplicate copy when the label is subjected to printing impact forces. However, in the Singer structure, care has to be taken to isolate the carbonless reactants from the pressure sensitive adhesive layer. Moreover, the labels of Singer have a tacky surface after delamination.\nLikewise, Lockhart, U.S. Pat. No. 4,277,089, teaches the use of a removable label product which is adhered to a backer sheet through the use of pressure sensitive adhesive. The backer sheet itself comprises a self-contained carbonless product which provides a duplicate copy when the face sheet is subjected to impact printing forces. Like Singer the peelable label of Lockhart has a tacky surface.\nAccordingly, the need still exists in the art for a clean release product which can provide a duplicate copy of information printed onto a face sheet or label and which can be readily removed from a laminate construction resulting in both a nontacky top sheet or label and a non-tacky web or backing sheet surface."}
{"text": "A MIMO (Multiple-Input, Multiple-Output) system is an example of a conventional communication system using multiple antennas. In multi-antenna communication, of which the MIMO system is typical, multiple transmission signals are each modulated, and each modulated signal is simultaneously transmitted from a different antenna in order to increase the transmission speed of the data.\nFIG. 23 illustrates a sample configuration of a transmission and reception device having two transmit antennas and two receive antennas, and using two transmit modulated signals (transmit streams). In the transmission device, encoded data are interleaved, the interleaved data are modulated, and frequency conversion and the like are performed to generate transmission signals, which are then transmitted from antennas. In this case, the scheme for simultaneously transmitting different modulated signals from different transmit antennas at the same time and on a common frequency is a spatial multiplexing MIMO system.\nIn this context, Patent Literature 1 suggests using a transmission device provided with a different interleaving pattern for each transmit antenna. That is, the transmission device from FIG. 23 should use two distinct interleaving patterns performed by two interleavers (πa and πb). As for the reception device, Non-Patent Literature 1 and Non-Patent Literature 2 describe improving reception quality by iteratively using soft values for the detection scheme (by the MIMO detector of FIG. 23).\nAs it happens, models of actual propagation environments in wireless communications include NLOS (Non Line-Of-Sight), typified by a Rayleigh fading environment is representative, and LOS (Line-Of-Sight), typified by a Rician fading environment. When the transmission device transmits a single modulated signal, and the reception device performs maximal ratio combination on the signals received by a plurality of antennas and then demodulates and decodes the resulting signals, excellent reception quality can be achieved in a LOS environment, in particular in an environment where the Rician factor is large. The Rician factor represents the received power of direct waves relative to the received power of scattered waves. However, depending on the transmission system (e.g., a spatial multiplexing MIMO system), a problem occurs in that the reception quality deteriorates as the Rician factor increases (see Non-Patent Literature 3).\nFIGS. 24A and 24B illustrate an example of simulation results of the BER (Bit Error Rate) characteristics (vertical axis: BER, horizontal axis: SNR (signal-to-noise ratio) for data encoded with LDPC (low-density parity-check) codes and transmitted over a 2×2 (two transmit antennas, two receive antennas) spatial multiplexing MIMO system in a Rayleigh fading environment and in a Rician fading environment with Rician factors of K=3, 10, and 16 dB. FIG. 24A gives the Max-Log approximation-based log-likelihood ratio (Max-log APP) BER characteristics without iterative detection (see Non-Patent Literature 1 and Non-Patent Literature 2), while FIG. 24B gives the Max-log APP BER characteristic with iterative detection (see Non-Patent Literature 1 and Non-Patent Literature 2) (number of iterations: five). FIGS. 24A and 24B clearly indicate that, regardless of whether or not iterative detection is performed, reception quality degrades in the spatial multiplexing MIMO system as the Rician factor increases. Thus, the problem of reception quality degradation upon stabilization of the propagation environment in the spatial multiplexing MIMO system, which does not occur in a conventional single-modulation signal system, is unique to the spatial multiplexing MIMO system.\nBroadcast or multicast communication is a service applied to various propagation environments. The radio wave propagation environment between the broadcaster and the receivers belonging to the users is often a LOS environment. When using a spatial multiplexing MIMO system having the above problem for broadcast or multicast communication, a situation may occur in which the received electric field strength is high at the reception device, but in which degradation in reception quality makes service reception difficult. In other words, in order to use a spatial multiplexing MIMO system in broadcast or multicast communication in both the NLOS environment and the LOS environment, a MIMO system that offers a certain degree of reception quality is desirable.\nNon-Patent Literature 8 describes a scheme for selecting a codebook used in precoding (i.e. a precoding matrix, also referred to as a precoding weight matrix) based on feedback information from a communication party. However, Non-Patent Literature 8 does not at all disclose a scheme for precoding in an environment in which feedback information cannot be acquired from the other party, such as in the above broadcast or multicast communication.\nOn the other hand, Non-Patent Literature 4 discloses a scheme for switching the precoding matrix over time. This scheme is applicable when no feedback information is available. Non-Patent Literature 4 discloses using a unitary matrix as the precoding matrix, and switching the unitary matrix at random, but does not at all disclose a scheme applicable to degradation of reception quality in the above-described LOS environment. Non-Patent Literature 4 simply recites hopping between precoding matrices at random. Obviously, Non-Patent Literature 4 makes no mention whatsoever of a precoding method, or a structure of a precoding matrix, for remedying degradation of reception quality in a LOS environment."}
{"text": "The present invention relates to a method of manufacturing color filters, and more particularly, to a method of increasing the adhesion of color filters and preventing cross talk effects on a semiconductor wafer.\nCharge-coupled device (CCDs) have been the mainstay of conventional imaging circuits for converting light into an electrical signal. The applications of CCDs include monitors, transcription machines and cameras. Although CCDs have many advantages, CCDs also suffer from high costs and the limitations imposed by its volume. To overcome the weakness of CCDs and reduce costs and dimensions, a CMOS photodiode device is developed. Since a CMOS photodiode device can be produced by using conventional techniques, both cost and the volume of the sensor can be reduced. The applications of CMOS photodiodes include PC cameras, digital cameras etc.\nWhether the image sensor device is composed of CCD or CMOS photodiode, the incident light must be separated into a combination of light of different wave lengths, for example, red, blue and green light. Then, the light is received by corresponding sensor devices and is transformed into electrical signals so as to obtain the original color of incident light by return of the electrical signals. Therefore, a color filter array must be formed on each photosensor device. Currently, color filters are produced by either patterning photosensitive resins using a photo-etching process with the resultant patterns being dyed by a dyeing material, or a photoresist containing dyeing material is directly used to produce color filters.\nPlease refer to FIG. 1 to FIG. 6. FIG. 1 to FIG. 6 are cross-sectional diagrams of manufacturing a color filter array on a photosensor device according to the prior art method. As shown in FIG. 1, the semiconductor wafer 10 contains a silicon substrate 12 and a P-well 14 positioned on the silicon substrate 12. The photosensor device contains a plurality of CMOS photodiodes and each photodiode contains a metal-oxide semiconductor (MOS) transistor (not shown) positioned on the P-well 14. A photosensor area 18 is formed on the P-well 14 to electrically connect with the MOS transistor. The MOS transistor is a complementary metal-oxide semiconductor (CMOS) transistor composed of an NMOS transistor and a PMOS transistor and functions as a CMOS transistor sensor. The semiconductor wafer 10 also contains a plurality of field oxide layers or shallow trench isolation (STI) structures 16 positioned on the silicon substrate 12 that surrounds the photo sensor area 18. The STI structures 16 act as a dielectric insulating material to prevent short circuiting due to contact between the photosensor areas 18 and other units.\nFirst, a passivation layer 20 is formed on the surface of the semiconductor wafer 10 that covers each photo sensor area 18. Next, as shown in FIG. 2, a red color filter layer (not shown) is formed on the surface of the semiconductor wafer 10. The color filter layer is composed of a positive type photoresist containing a red dye in a large amount (dry weight) of 10 to 50 wt %. A pattern-exposure process is used to form patterns of red color filters in the color filter layer, then the exposed portions of the filter layer is removed to form each red color filter 22. For increasing the effect and reliability of color filters, an ultraviolet (UV) light irradiation and heating process is performed after the formation of the red color filters 22. The UV light used has a wavelength of 320 nm or less at a quantity of 20 J/cm2 or less. The heating process is preferably performed in an inert atmosphere, for example, in nitrogen (N2) for suppressing the oxidation of the photoresist material. The starting temperature of the heating process is between a range of 60xc2x0 C. to 140xc2x0 C. Then, an average increasing temperature rate used in the heating process is 1.5xc2x0 C./sec. The end temperature of the heating process is between a range of 160xc2x0 C. to 220xc2x0 C.\nNext, green and blue color filters are formed by repeating the above-mentioned processes. As shown in FIG. 3, a green color filter layer 24 is formed on the surface of the semiconductor wafer 10. The color filter layer 24 is composed of a positive type photoresist containing a green dye in a large amount (dry weight) of 10 to 50 wt %. As shown in FIG. 4, a pattern-exposure process is used to form patterns of green color filters in the color filter layer 24, then the exposed portions of the filter layer 24 is removed to form each green color filter 26. For increasing the effect and reliability of color filters, a UV light irradiation and heating process is also performed after the formation of green color filters 26. The UV light used has a wavelength of 320 nm or less at a quantity of 20 J/cm2 or less. The heating process is preferably performed in an inert atmosphere, for example, in nitrogen (N2) for suppressing the oxidation of the photoresist material. The starting temperature in heating process is between a range of 60xc2x0 C. to 140xc2x0 C. Then, an average increasing temperature rate used in the heating process is 1.5xc2x0 C./sec. The end temperature of the heating process is between a range of 160xc2x0 C. to 220xc2x0 C.\nAs shown in FIG. 5, a blue color filter layer 28 is formed on the surface of the semiconductor wafer 10. The color filter layer 28 is composed of a positive type photoresist containing a blue dye in a large amount (dry weight) of 10 to 50 wt %. As shown in FIG. 6, a pattern-exposure process is used to form patterns of blue color filters in the color filter layer 28, then the exposed portions of the filter layer 28 is removed to form each blue color filter 30. For increasing the effect and reliability of color filters, a UV light irradiation and heating process is also performed after the formation of the blue color filters 30. The UV light used for irradiation has a wavelength of 320 nm or less at a quantity of 20 J/cm2 or less. The heating process is preferably performed in an inert atmosphere, for example, in nitrogen (N2) for suppressing the oxidation of the photoresist material. The starting temperature of the heating process is between a range of 60xc2x0 C. to 140xc2x0 C. Then, an average increasing temperature rate used in the heating process is 1.5xc2x0 C./sec. The end temperature of the heating process is between a range of 160xc2x0 C. to 220xc2x0 C. The color filter array of a photosensor device produced by the prior art method is then completed.\nColor filters produced by the prior art method have only the bottom surfaces contacting the passivation layer. Therefore, adhesion is weak and stripping of the color filters easily occurs. Also, a space exists between each color filter, so that scattered light easily penetrates through the space to enter the neighboring photosensor area and result in cross talk effects. The noise signals received by CMOS transistor sensors increase and the sensitivity is reduced.\nIt is therefore a primary objective of the present invention to provide a method of manufacturing color filters for increasing adhesion of color filters and preventing stripping of the color filters. A barrier layer is simultaneously formed to prevent incident light from penetrating the space between the color filters to cause cross talk effects.\nThe present invention provides a method for increasing both the adhesion of color filters and preventing cross talk effects on a semiconductor wafer. The semiconductor wafer comprises a substrate, a plurality of metal-oxide semiconductor (MOS) transistor sensors positioned on the substrate, and a plurality of insulators positioned on the substrate. Each insulator is positioned between two MOS transistor sensors. The method first involves forming a dielectric layer on the semiconductor wafer, which covers each MOS transistor sensor and insulators. Then, a plurality of metal layers are formed on the dielectric layer, each two metal layers positioned approximately above two ends of one MOS transistor sensor. Next, a passivation layer is formed, followed by the formation of a first color filter on the passivation layer and above one MOS transistor sensor. Two ends of the first color filter are aligned approximately above the two metal layers. Thereafter, a second color filter and a third color filter are sequentially formed on the passivation layer, and portions of both the second color filter and the third color filter cover one end of the first color filter. The portion of the second color filter covering the first color filter and the portion of the third color filter covering the first color filter are both used to enhance adhesion of the color filters, and simultaneously, the portion of the second color filter covering the first color filter and the portion of the third color filter covering the first color filter are both used as barriers to prevent cross talk effects.\nThe color filters produced by the present invention have portions covering one another so as to increase both the contact areas and the adhesion of the color filters. As well, the covered portions of each color filter are used as a barrier layer to prevent cross talk effects.\nThese and other objectives of the present invention will no doubt become obvious to those of ordinary skill in the art after reading the following detailed description of the preferred embodiment, which is illustrated in the various figures and drawings."}
{"text": "Aspects of the present disclosure relate generally to optical fiber cables, interconnect assemblies including the optical fiber cables, and methods of manufacturing the interconnect assemblies.\nInterconnect cable assemblies are typically used to connect servers and other hardware in data centers and other facilities. Due to industry trends of increasing speed, memory, and density of hardware in data centers, management and organization of correspondingly increasing numbers of optical fiber cables as well as complexity of the associated interconnections becomes more and more of an issue. Accordingly interconnect cable assemblies have been developed to help control and organize the optical fiber cables and interconnections in a data center. One such development includes harness cable assemblies.\nA harness cable assembly is a particular type of interconnect cable assembly where one end of the harness cable assembly includes a multi-fiber connector that supports optical fibers extending through the harness cable assembly to the other end, which includes furcated legs each having local connectors that support subsets of the optical fibers. A jacket holds discrete subsets of the optical fibers together in modules so that the optical fibers may be efficiently routed together through the data center. Routing the cables together helps to minimize the volume of space used by the optical fiber cables and correspondingly helps to organize densely populated data centers and other facilities. However, conventional manufacturing of harness cable assemblies is a time-consuming and resource-intensive process, which may suffer from low process yields.\nConventionally to manufacture a harness cable assembly, manufacturers first select a length of optical fiber cable. The cable typically includes a jacket surrounding buffer tubes that each include optical fibers. Within the jacket, the buffer tubes are tightly packed together to allow for a minimal jacket diameter, saving space and material costs. The buffer tubes of optical fiber cables are typically stranded. The manufacturers attach local connectors to the optical fibers of the buffer tubes, typically removing about three to seven meters (e.g., ten to twenty feet) of the cable jacket and feeding the optical fibers through an empty furcation tube. However, this approach only allows for a limited length of the harness assembly. In some such cases, the maximum length allowed by this method is about eight meters (e.g., twenty-five feet), because beyond about eight meters, friction generated by pushing the optical fibers through the furcation tube may break the fibers. Furthermore, the lengths of the furcated legs are not adjustable or customizable with the conventional method once local connectors are attached to the legs—thus ruining desired configurations of the legs, such as a precise staggering formations, if a local connector on any of the legs needs to be reworked (e.g., replaced, reattached, etc.) during manufacturing.\nAccordingly, a need exists for an interconnect optical fiber cable assembly structure or manufacturing method that allows for longer cable assembly lengths. A need exists for an interconnect optical fiber cable assembly structure or manufacturing method that allows for adjustment of the leg lengths once local connectors are attached to the legs, if a local connector on any of the legs needs to be reworked during manufacturing. A need exists for an interconnect optical fiber cable assembly structure or manufacturing method that reduces time and manufacturing costs of producing the interconnect optical fiber cable, and thereby increasing process yields."}
{"text": "The present invention relates to new multiheterocyclic compounds as integrin antagonists with a broad spectrum of action having, inter alia, antiosteoporotic, antirestenotic, anticarcinogenic and antiatherosclerotic activity. The present invention moreover relates to the preparation of these compounds and their use for the production of medicaments, and also medicaments comprising them.\nIntegrins are heterodimeric transmembrane proteins found on the surface of cells, which play an important part in the adhesion of the cells to an extracellular matrix. They recognize extracellular glycoproteins such as fibronectin or vitronectin on the extracellular matrix by means of the RGD sequence occurring in these proteins (RGD is the single letter code for the amino acid sequence arginine-glycine-aspartate).\nIn general, integrins such as, for example, the vitronectin receptor, which is also called the xcex1vxcex23 receptor, or alternatively the xcex1vxcex25 receptor or the GpIIb/IIIa receptor play an important part in biological processes such as cell migration and cell-matrix adhesion and thus in diseases in which these processes are crucial steps. Cancer, osteoporosis, arteriosclerosis, restenosis (reoccurrence of stenosis after percutaneous transluminal angioplasty) and opthalmia may be.mentioned by way of example.\nThe xcex1vxcex23 receptor occurs, for example, in large amounts on growing endothelial cells and makes possible their adhesion to an extracellular matrix. Thus the xcex1vxcex23 receptor plays an important part in angiogenesis, i.e. the formation of new blood vessels, which is a crucial prerequisite for tumor growth and metastasis formation in carcinoses. Furthermore, it is also responsible for the interaction between osteoclasts, i.e. cells resorbing mineralized tissue, and the bone structure. The first step in the degradation of bone tissue consists in the adhesion of osteoclasts to the bone. This cell-matrix interaction takes place via the xcex1vxcex23 receptor, which is why the corresponding integrin plays an important part in this process. Osteolytic diseases such as osteoporosis are induced by an inequilibrium between bone formation and bone destruction, i.e. the resorption of bone material caused by accumulation of osteoclasts predominates.\nIt was possible to show that the blockage of the abovementioned receptors is an important starting point for the treatment of disorders of this type. If the adhesion of growing endothelial cells to an extracellular matrix is suppressed by blocking their appropriate integrin receptors, for example, by a cyclic peptide or a monoclonal antibody, the endothelial cells die. Therefore angiogenesis does not occur, which leads to a cessation or resolution of the tumor growth (cf., for example, Brooks et al., Cell, Volume 79, 1157-1164, 1994).\nMoreover, the invasive properties of tumor cells and thus their capability for metastasis formation are markedly decreased if their xcex1vxcex23 receptor is blocked by an antibody (Brooks et al., J. Clin. Invest., Volume 96, 1815, 1995).\nThe degradation of bone tissue can be suppressed by blockage of the xcex1vxcex23 receptors of the osteoclasts, since these are then unable to accumulate on the bone in order to absorb its substance (WO 98/18461, p. 1, 1.24 to p. 2, 1.13).\nBy means of the blockage of the xcex1vxcex23 receptor on cells of the smooth aorta vascular musculature with the aid of integrin receptor antagonists, the migration of these cells into the neointima and thus angioplasty leading to arteriosclerosis and restenosis can be suppressed (Brown et al., Cardiovascular Res., Volume 28, 1815, 1994).\nIn recent years, compounds have therefore been sought which act as antagonists of integrin receptors. For example WO 97/24119 discloses victronectin (xcex1vxcex23) receptor antagonist as antiosteoporosis agents which are structurally related to the compounds described here, but do not have a broad spectrum of action of this type.\nIt was the object of the present invention to develop compounds which exhibit a high activity as integrin antagonists and in particular against the xcex1vxcex23 and/or the xcex1vxcex25 receptor.\nThe compounds on which the present invention is based can be described by the following general formula (I): \nwhere\nR1 and R2 together with the formal double bond bridging them form a phenyl ring or a six-membered aromatic heterocycle having 1 or 2 nitrogen atoms, which is optionally substituted by halogen, nitro, cyano, trifluoromethyl, hydroxyl, (C1-C6)-alkyl, (C1-C6)-alkoxy, (C1-C6)-alkoxycarbonyl, amino, carboxyl, phenoxy, (C6-C10)-aryl, (C1-C6)-alkylamino, (C1-C6)-alkylsulfone or a group of the formula xe2x80x94SO2NRaRb wherein Ra and Rb may independently represent hydrogen or (C1-C6)-alkyl, and\nR3=H or (C1-C4)-alkyl, and\nA=O, S, (CH2)n where n=1,2,3 or 4 or Nxe2x80x94R7 where R7=H or (C1-C4)-alkyl, or is absent and\nX, Y, Z=O, S, N, Nxe2x80x94R8 where R8=H, (C1-C6)-alkyl or (C6-C10)-aryl, or Cxe2x80x94(R9)(R10) where R9, R10=H, (C1-C6)-alkyl or (C1-C10)-aryl, or Cxe2x80x94(R11) if the X, Y or Z in question is participating in a formal double bond, where R11=H, (C1-C6)-alkyl or (C6-C10)-aryl, and\nD=O, S, (CH2)n where n=1,2,3 or 4, or Nxe2x80x94R7 where R7=H or (C1-C4)-alkyl, or is absent and\nE, G, L=H, halogen, nitro, cyano, trifluoromethyl, hydroxyl, carboxyl, (C1-C6)-alkyl, (C1-C6)-alkoxy or (C1-C6)-alkoxycarbonyl, where E, G and L can be identical or different and\nR4=H, (C3-C8)-cycloalkyl, (C1-C6)-alkyl, which is optionally substituted by hydroxyl, (C1-C6)-alkoxy or phenyl, where the latter is optionally in turn substituted on the phenyl ring by halogen, (C1-C4)-alkyl, (C1-C4)-alkoxy, hydroxyl, trifluoromethoxy or trifluoromethyl, and\nR5=H or (C1-C4)-alkyl and\nR6=H, (C1-C6)-alkyl or benzyl\nand their salts.\nPreferred compounds according to the general formula (I) are those where\nR1 and R2 together with the formal double bond bridging them form a phenyl ring or a six-membered aromatic heterocycle having 1 or 2 nitrogen atoms, which is optionally substituted by (C1-C6)-alkoxy, nitro or amino, and\nR3=H, and\nA=O, S, (CH2)n where n=1,2,3 or 4 or Nxe2x80x94H, or is absent and\nX, Y, Z=O, S, N, Nxe2x80x94R8 where R8=H, (C1-C6)-alkyl or (C6-C10)-aryl, or Cxe2x80x94(R9)(R10) where R9, R10=H, (C1-C6)-alkyl or (C6-C10)-aryl, or Cxe2x80x94(R11) if the X, Y or Z in question is participating in a formal double bond, where R11=H, (C1-C6)-alkyl or (C6-C10)-aryl, and\nD=O, S, (CH2)n where n=1, 2, 3 or 4, or Nxe2x80x94H, or is absent and\nE, G, L=H, and\nR4=H, benzyl, cyclopropyl, cyclopentyl, cyclohexyl, phenylethyl or (C1-C4)-alkyl, which is optionally substituted by hydroxyl or methoxy, and\nR5=H, and\nR6=H or methyl.\nand their salts.\nParticularly preferred compounds according to the general formula (I) are those where\nR1 and R2 together with the formal double bond bridging them form a phenyl ring, which is optionally substituted by nitro. and\nR3=H and\nA=is absent and\nX=N, Y=O and Z=N or X=N, Y=N and Z=O, S or X=CH2, Y=O, Z=N and\nD=is absent and\nE, G, L=H, and\nR4=methyl, cyclopropyl or (C1-C4)-alkyl substituted by methoxy and\nR5=H, and\nR6=H or methyl\nand their salts.\nThe compounds according to the invention can also be present in the form of their salts. In general, salts with organic or inorganic bases or acids may be mentioned here.\nIn the context of the present invention, physiologically acceptable salts are preferred. Physiologically acceptable salts of the compounds according to the invention can be salts of the substances according to the invention with mineral acids, carboxylic acids or sulphonic acids. Particularly preferred salts are, for example, those with hydrochloric acid, hydrobromic acid, sulphuric acid, phosphoric acid, methanesulphonic acid, ethanesulphonic acid, toluenesulphonic acid, benzenesulphonic acid, naphthalenedisulphonic acid, acetic acid, propionic acid, lactic acid, tartaric acid, citric acid, fumaric acid, maleic acid or benzoic acid.\nPhysiologically acceptable salts can also be metal or ammonium salts of the compounds according to the invention. Those particularly preferred are, for example, sodium, potassium, magnesium or calcium salts, and also ammonium salts which are derived from ammonia, or organic amines, such as, for example, ethylamine, di- or triethylamine, di- or triethanolamine, dicyclohexylamine, dimethylaminoethanol, arginine, lysine, ethylenediamine or 2-phenylethylamine.\nThe present invention also includes ammonium compounds which can be prepared by conversion of the free amines by means of alkylation.\nIn the context of the present invention, the substituents in general have the following meaning:\n(C1-C6)-alkyl in general represents, depending on the abovementioned substituents, a straight-chain or branched hydrocarbon radical having 1 to 6, preferably 1 to 4 carbon atoms. Examples which may be mentioned are methyl, ethyl, propyl isopropyl, butyl, isobutyl, pentyl, isopentyl, hexyl, isohexyl.\n(C1-C6)-alkoxycarbonyl can be represented, for example, by the formula \nAlkyl here represents a straight-chain or branched hydrocarbon radical having 1 to 6 carbon atoms. Lower alkoxycarbonyl having 1 to 4 carbon atoms in the alkyl moiety is preferred. Examples which may be mentioned are the following alkoxycarbonyl radicals: methoxycarbonyl, ethoxycarbonyl, propoxycarbonyl, isopropoxycarbonyl, butoxycarbonyl or isobutoxycarbonyl.\n(C3-C8)-cycloalkyl in general represents a cyclic hydrocarbon radical having 3 to 8 carbon atoms. Cyclopropyl, cyclopentyl and cyclohexyl are preferred. Examples which may be mentioned are cyclopentyl, cyclohexyl, cycloheptyl and cyclooctyl.\n(C6-C10)-aryl in general represents an aromatic radical having 6 to 10 carbon atoms. Preferred aryl radicals are phenyl and naphthyl.\nHalogen in the context of the invention represents fluorine, chlorine, bromine and iodine.\n(C1-C6)-alkoxy in general represents, depending on the abovementioned substituents, a straight-chain or branched hydrocarbon radical bonded via an oxygen atom and having 1 to 6 carbon atoms. Lower alkoxy having 1 to 4 carbon atoms is preferred. An alkoxy radical having 1 to 3 carbon atoms is particularly preferred. Examples which may be mentioned are methoxy, ethoxy, propoxy, isopropoxy, butoxy, isobutoxy, pentoxy, isopentoxy, hexoxy, isohexoxy.\n(C1-C6)-alkylsulfone can be represented by the formula \nwherein Alkyl represents a straight-chain or branched hydrocarbon radical having 1 to 6, preferably 1 to 4 carbon atoms. Examples which may be mentioned are: methylsulfone, ethylsulfone, propylsulfone, isopropylsulfone.\nSix-membered aromatic heterocycles in the context of the invention, depending on the abovementioned substituents, in general represent a 6-membered heterocycle which can contain up to 2 nitrogen atoms and which can optionally also be bonded via a nitrogen atom. Examples which may be mentioned are: pyridyl and pyrimidyl.\nIn general formula (I) the heterocyclic moiety \nrepresents an aromatic or partially saturated heterocycle. Preferred examples for this moiety are: \nThe compounds according to the invention can exist in stereoisomeric forms which either behave as image and mirror image (enantiomers), or which do not behave as image and mirror image (diastereomers). The invention relates both to the enantiomers or diastereomers and their respective mixtures. These mixtures of the enantiomers and diastereomers can be separated into the stereoisomerically uniform constituents in a known manner.\nThe compounds of the general formula (I) according to the invention can be prepared in the following way using process variant (A) or (B):\n(A) reaction of a component (a) \nxe2x80x83with a component (b) \nwhere (b) is first reacted in the presence of sulfonyl chloride, POCl3 or PCl5 to give the acid chloride and is then reacted with (a) in the presence of a base in an inert organic solvent, or a coupling of (a) with (b) is carried out by means of Bis(2-oxo-3-oxazolidinyl)phosphinic chloride (BOP) or carbodiimide reagents, followed by a cyclization reaction which can be effected by means of heat, acids, bases, dehydrating substances and by addition of generally nucleophilic reagents,\nwhere, if appropriate, the COOR6 group is finally hydrolyzed and where Q in one of the two components (a) or (b) is oxygen and in the other components corresponds to Z, or\n(B) reaction of a component (axe2x80x2) \nxe2x80x83with a component (b) \nwhere (b) is first reacted in the presence of sulfonyl chloride, POCl3 or PCl5 to give the acid chloride and then reacted with (axe2x80x2) in the presence of a base in an inert organic solvent, or a coupling of (axe2x80x2) with (b) is carried out by means of BOP or carbodiimide reagents, followed by a cyclization reaction which can be effected by means of heat, acids, bases, dehydrating substances and by addition of generally nucleophilic reagents,\nwhere, if appropriate, the COOR6 group is finally hydrolyzed and where Q in one of the two components (axe2x80x2) or (b) is oxygen and in the other component corresponds to Z.\nPreferably, in process variant (A), for component (a), Xxe2x95x90N, Yxe2x95x90O and Qxe2x95x90NH is employed and for component (b), Qxe2x95x90O is employed and in process variant (B), for component (axe2x80x2), Xxe2x95x90NH, Yxe2x95x90NH and Qxe2x95x90O is employed and, for component (bxe2x80x2), Qxe2x95x90O is employed.\nParticularly preferably, it applies to both process variants that R1 and R2 together with the formal double bond bridging them form a phenyl radical,\nR3, R5, E, G, L=H,\nR4, R6 =methyl and\nA, D=is absent.\nThe compounds of the formula (I) according to the invention have a surprisingly wide spectrum of pharmacological action.\nThe compounds according to the invention exhibit an antagonistic action against integrin inhibitors, in particular the xcex1vxcex23 receptor or the xcex1vxcex25 receptor.\nThey can be used as active compounds in medicaments for the reduction of changes to vascular walls and are employed for the treatment of arterial hypertension and atherosclerosis. Moreover, they can be employed for the treatment of coronary heart disorders, cardiac insufficiency, disorders of brain function, ischemic brain disorders, (peripheral) circulation disorders, microcirculation disorders and thromboses, functional disorders of the kidney and adrenal gland, bronchospastic and vascular system-related disorders of the airways, sodium retention and edemas as well as osteolytic disorders such as osteoporosis, cancer, carcinoses and ophthalmic diseases.\nFurthermore, the proliferation and migration of smooth muscle cells plays a decisive part in the occlusion of vessels. The compounds according to the invention are suitable for inhibiting this proliferation and can therefore also be employed for the treatment of restenosis.\nThe novel active compounds are distinguished pharmacologically by good kinetic parameters. In particular, they have favorable properties with respect to clearance.\nThe novel active compounds can be converted in a known manner into the customary formulations, such as tablets, coated tablets, pills, granules, aerosols, syrups, emulsions, suspensions and solutions, using inert, nontoxic, pharmaceutically suitable excipients or solvents. In this case the therapeutically active compound should in each case be present in a concentration of approximately 0.5 to 90% by weight of the total mixture, i.e. in amounts which are sufficient in order to achieve the dosage range indicated.\nThe formulations are prepared, for example, by extending the active compounds with solvents and/or excipients, if appropriate using emulsifiers and/or dispersants, where, for example, if water is used as a diluent organic solvents can optionally be used as auxiliary solvents.\nAdministration is carried out in a customary manner, preferably orally or parenterally, in particular perlingually or intravenously.\nIn the case of parenteral administration, solutions of the active compound using suitable liquid carrier materials can be employed.\nIn general, it has proven advantageous in the case of intravenous administration to administer amounts from approximately 0.001 to 1 mg/kg, preferably approximately 0.01 to 0.5 mg/kg, of body weight to achieve effective results, and in the case of oral administration the dose is approximately 0.01 to 100 mg/kg, preferably 0.1 to 50 mg/kg, of body weight.\nIn spite of this, if appropriate, it may be necessary to depart from the amounts mentioned, namely depending on the body weight or the type of application route, on individual behavior toward the medicament, the manner of its formulation and the time or interval at which administration takes place. Thus in some cases it may be adequate to manage with less than the abovementioned minimum amount, while in other cases the upper limits mentioned must be exceeded. In the case of the administration of relatively large amounts, it may be advisable to divide these into several individual doses over the course of the day.\nThe present invention is illustrated in greater detail below by working examples.\nStarting Compound I \n0.877 g of methyl S-7-carboxy-2,3,4,5-tetrahydro-4-methyl-3-oxo-1H-1,4-benzodiazepine-2-acetate (Tetrahedron Letters 38, 3131-34, 1997) is stirred under reflux for 1 h with thionyl chloride (40 ml), concentrated and the residue is dissolved in 20 ml of chloroform. The solution obtained is added dropwise to a mixture of 0.53 g of benzimidazole-2-carboxamidoxime (J. Chem. Soc. C, 1967, 28) and chloroform (50 ml), pyridine (0.48 g) and triethylamine (0.30 g). It is stirred at room temperature for 2 h, then concentrated and the residue is taken up in ethyl acetate. After this, washing with water, saturated sodium hydrogencarbonate solution and saturated sodium chloride solution and drying over magnesium sulfate takes place, after concentration an amorphous residue being obtained which is stirred with dichloromethane and filtered off with suction. Recrystallization is carried out from ethanol (0.41 g). The mother liquor is concentrated and chromatographically purified (CH2Cl2; CH2 Cl2/CH3OH=30/1). 0.6 g of the title compound is obtained from the mother liquor."}
{"text": "Many integrated circuits include a test access port that may be used to verify the proper operation of the integrated circuit and the proper operation of the system that contains the integrated circuit, including the printed circuit board on which the integrated circuit is mounted. A test access port provides a test clock (TCK) signal, a test mode select (TMS) signal, a test data input (TDI) signal, and a test data output (TDO) signal. The TDI signal may be used to serially shift a bit of data into the integrated circuit each rising edge of the TCK signal and the TDO signal may be used to serially shift a bit of data out of the integrated circuit each falling edge of the TCK signal. A test access port may be implemented according to the IEEE 1149.1 standard and derived standards, such as IEEE 1532 and IEEE 1149.6.\nAn operating mode for the test access port may determine the destination for data shifted into the integrated circuit at the TDI signal and the source for data shifted out of the integrated circuit at the TDO signal. The operating mode of the test access port may be specified by one of sixteen states of a state vector of a finite state machine. Each state of the finite state machine transitions to one of two possible next states on each rising edge of the TCK signal, with the next state for a current state of the finite state machine determined by the value of the TMS signal at the rising edge of the TCK signal. Thus, the TMS signal may be used to control the operating mode of the test access port.\nWhile the test access port is very useful for determining the cause of improper operation for many aspects of an integrated circuit, the test access port may have limited usefulness for determining the cause of improper operation related to the test access port itself. For example, electrical noise may cause the finite state machine of the test access port to make improper transitions between states, such that the operating mode of the test access port cannot be properly controlled by the TMS signal. Determining the cause of improper operation may be time-consuming and expensive when the operating mode of the test access port cannot be properly controlled due to such electrical noise. The present invention may address one or more of the above issues."}
{"text": "This invention relates to a drive mechanism for at least two windshield wipers of a vehicle, in particular of an automotive vehicle.\nFrom DE-GM 7103702, a drive mechanism for two windshield wipers is known which includes a driving motor with a gear unit. The gear unit is accommodated in a gear casing which is directly linked to the driving motor and which further is designed so that two wiper bearings are worked into the gear casing and the gear casing also receives the driving rods for the two wiper bearings. The gear casing which is a casting or molded part is entirely closed by a cap so that, on principle, the overall drive mechanism features a completely molded case design. Due to its compact design, the essential advantage of such a drive mechanism is the easy handling when being mounted in a vehicle. On the other hand, the gear casing features a relatively complicated shape so that the costs of its manufacture are relatively high. Moreover, the shape of the gear casing invariably determines the points of support of the wiper shafts or rather the distance between the wiper shafts. Thus, a change in the points of support of the wiper shafts or rather in the geometry of the drive mechanism always implies the manufacture of a gear casing of different shape. This known drive mechanism, thus, cannot easily be adapted to varying geometrical conditions of different vehicles so that this drive mechanism, on principle, can be used for the same vehicle, only.\nAnother drive mechanism for at least two windshield wipers is known from EP 0 409 944 B1. This drive mechanism also includes a driving motor comprising a gear unit which is accommodated in a gear casing which is a casting or molded part and is closed by a cap. The drive mechanism includes two wiper bearings each of which is linked to respectively the gear casing or the cap of the gear casing via a tubular carrier. Each of the tubular carriers is slipped with its opposite ends onto a pivot connected with the wiper bearing and onto a pivot connected with respectively the gear casing or the gear casing cap and is fixed thereon by means of caulking. The essential advantage of this drive mechanism is that it is adaptable in an extremely flexible manner to very different geometrical conditions of different vehicles. On the one hand, such adaptations to the space conditions in the respective vehicle are made by appropriately selecting the length of the carriers which connect the wiper bearings with respectively the gear casing or the gear casing cap. On the other hand, it is possible to effect a further adaptation by bending the carriers or by changing the angle position of the wiper bearings around the longitudinal axis of the carrier. A disadvantage of this drive mechanism is the relatively big number of individual component parts from which this mechanism is assembled. This relatively big number of individual component parts involves increased assembly efforts for manufacturing this mechanism, a heavier weight of the drive mechanism as well as a larger amount of arrangements as regards the individual component parts.\nIt now is an object of this invention to improve a drive mechanism for at least two windshield wipers so that, on the one hand, it will be relatively easy to adapt it to the space conditions in different vehicles and, on the other hand, to ensure a low-cost manufactureabilty by reducing the required number of component parts."}
{"text": "Polymers, especially synthetic plastics, are ubiquitous in daily life due to their relatively low production costs and good balance of material properties. Synthetic plastics are used in a wide variety of applications, such as packaging, automotive components, medical devices, and consumer goods. To meet the high demand of these applications, tens of billions of pounds of synthetic plastics are produced globally on an annual basis. The overwhelming majority of synthetic plastics are produced from increasingly scarce fossil sources, such as petroleum and natural gas. Additionally, the manufacturing of synthetic plastics from fossil sources produces CO2 as a by-product.\nThe ubiquitous use of synthetic plastics has consequently resulted in millions of tons of plastic waste being generated every year. While the majority of plastic waste is landfilled via municipal solid waste programs, a significant portion of plastic waste is found in the environment as litter, which is unsightly and potentially harmful to ecosystems. Plastic waste is often washed into river systems and ultimately out to sea.\nPlastics recycling has emerged as one solution to mitigate the issues associated with the wide-spread usage of plastics. Recovering and re-using plastics diverts waste from landfills and reduces the demand for virgin plastics made from fossil-based resources, which consequently reduces greenhouse gas emissions. In developed regions, such as the United States and the European Union, rates of plastics recycling are increasing due to greater awareness by consumers, businesses, and industrial manufacturing operations. The majority of recycled materials, including plastics, are mixed into a single stream which is collected and processed by a material recovery facility (MRF). At the MRF, materials are sorted, washed, and packaged for resale. Plastics can be sorted into individual materials, such as high-density polyethylene (HDPE) or poly(ethylene terephthalate) (PET), or mixed streams of other common plastics, such as polypropylene (PP), low-density polyethylene (LDPE), poly(vinyl chloride) (PVC), polystyrene (PS), polycarbonate (PC), and polyamides (PA). The single or mixed streams can then be further sorted, washed, and reprocessed into a pellet that is suitable for re-use in plastics processing, for example blow and injection molding.\nThough recycled plastics are sorted into predominately uniform streams and are washed with aqueous and/or caustic solutions, the final reprocessed pellet often remains highly contaminated with unwanted waste impurities, such as spoiled food residue and residual perfume components. In addition, recycled plastic pellets, except for those from recycled beverage containers, are darkly colored due to the mixture of dyes and pigments commonly used to colorize plastic articles. While there are some applications that are insensitive to color and contamination (for example black plastic paint containers and concealed automotive components), the majority of applications require non-colored pellets. The need for high quality, “virgin-like” recycled resin is especially important for food and drug contact applications, such as food packaging. In addition to being contaminated with impurities and mixed colorants, many recycled resin products are often heterogeneous in chemical composition and may contain a significant amount of polymeric contamination, such as polyethylene (PE) contamination in recycled PP and vice versa.\nMechanical recycling, also known as secondary recycling, is the process of converting recycled plastic waste into a re-usable form for subsequent manufacturing. A more detailed review of mechanical recycling and other plastics recovery processes are described in S. M. Al-Salem, P. Lettieri, J. Baeyens, “Recycling and recovery routes of plastic solid waste (PSW): A review”, Waste Management, Volume 29, Issue 10, October 2009, Pages 2625-2643, ISSN 0956-053X. While advances in mechanical recycling technology have improved the quality of recycled polymers to some degree, there are fundamental limitations of mechanical decontamination approaches, such as the physical entrapment of pigments within a polymer matrix. Thus, even with the improvements in mechanical recycling technology, the dark color and high levels of chemical contamination in currently available recycled plastic waste prevents broader usage of recycled resins by the plastics industry.\nTo overcome the fundamental limitations of mechanical recycling, there have been many methods developed to purify contaminated polymers via chemical approaches, or chemical recycling. Most of these methods use solvents to decontaminate and purify polymers. The use of solvents enables the extraction of impurities and the dissolution of polymers, which further enables alternative separation technologies.\nFor example, U.S. Pat. No. 7,935,736 describes a method for recycling polyester from polyester-containing waste using a solvent to dissolve the polyester prior to cleaning. The '736 patent also describes the need to use a precipitant to recover the polyester from the solvent.\nIn another example, U.S. Pat. No. 6,555,588 describes a method to produce a polypropylene blend from a plastic mixture comprised of other polymers. The '588 patent describes the extraction of contaminants from a polymer at a temperature below the dissolution temperature of the polymer in the selected solvent, such as hexane, for a specified residence period. The '588 patent further describes increasing the temperature of the solvent (or a second solvent) to dissolve the polymer prior to filtration. The '588 patent yet further describes the use of shearing or flow to precipitate polypropylene from solution. The polypropylene blend described in the '588 patent contained polyethylene contamination up to 5.6 wt %.\nIn another example, European Patent Application No. 849,312 (translated from German to English) describes a process to obtain purified polyolefins from a polyolefin-containing plastic mixture or a polyolefin-containing waste. The '312 patent application describes the extraction of polyolefin mixtures or wastes with a hydrocarbon fraction of gasoline or diesel fuel with a boiling point above 90° C. at temperatures between 90° C. and the boiling point of the hydrocarbon solvent. The '312 patent application further describes contacting a hot polyolefin solution with bleaching clay and/or activated carbon to remove foreign components from the solution. The '312 patent yet further describes cooling the solution to temperatures below 70° C. to crystallize the polyolefin and then removing adhering solvent by heating the polyolefin above the melting point of the polyolefin, or evaporating the adhering solvent in a vacuum or passing a gas stream through the polyolefin precipitate, and/or extraction of the solvent with an alcohol or ketone that boils below the melting point of the polyolefin.\nIn another example, U.S. Pat. No. 5,198,471 describes a method for separating polymers from a physically commingled solid mixture (for example waste plastics) containing a plurality of polymers using a solvent at a first lower temperature to form a first single phase solution and a remaining solid component. The '471 patent further describes heating the solvent to higher temperatures to dissolve additional polymers that were not solubilized at the first lower temperature. The '471 patent describes filtration of insoluble polymer components.\nIn another example, U.S. Pat. No. 5,233,021 describes a method of extracting pure polymeric components from a multi-component structure (for example waste carpeting) by dissolving each component at an appropriate temperature and pressure in a supercritical fluid and then varying the temperature and/or pressure to extract particular components in sequence. However, similar to the '471 patent, the '021 patent only describes filtration of undissolved components.\nIn another example, U.S. Pat. No. 5,739,270 describes a method and apparatus for continuously separating a polymer component of a plastic from contaminants and other components of the plastic using a co-solvent and a working fluid. The co-solvent at least partially dissolves the polymer and the second fluid (that is in a liquid, critical, or supercritical state) solubilizes components from the polymer and precipitates some of the dissolved polymer from the co-solvent. The '270 patent further describes the step of filtering the thermoplastic-co-solvent (with or without the working fluid) to remove particulate contaminants, such as glass particles.\nThe known solvent-based methods to purify contaminated polymers, as described above, do not produce “virgin-like” polymer. In the previous methods, co-dissolution and thus cross contamination of other polymers often occurs. If adsorbent is used, a filtration and/or centrifugation step is often employed to remove the used adsorbent from solution. In addition, isolation processes to remove solvent, such as heating, vacuum evaporation, and/or precipitation using a precipitating chemical are used to produce a polymer free of residual solvent.\nAccordingly, a need still exists for an improved solvent-based method to purify contaminated polymers that uses a solvent that is readily and economically removed from the polymer, is relatively simple in terms of the number of unit operations, produces a polymer without a significant amount of polymeric cross contamination, produces a polymer with improved color properties, and produces a polymer with reduced malodor."}
{"text": "Typical machining operations such as milling and routing utilize specially designed tool bits that have a fixed radii cutting feature thereon. A problem with such traditional tools is that most machined parts have varying contours such as a changing radius on the edge concave, a fillet at a corner, multiple radii and other rapidly changing features which require a machinist to find and then to frequently exchange the tool bit with an appropriate one to form the final shape of the part. This is time consuming and frustrating particularly when the machinist must search through dozens of tool bits until the correct one is located. Furthermore, if a particular tool bit is unsuitable for the cutting task, the machinist oftentimes grinds and machines the tool bit in an attempt to simulate the desired radii to be cut. This, too, is time consuming and ruins the tool bit as well as compromising the accuracy of the workpiece.\nNo device is known such as a variable radii cutter that conveniently allows cutting a wide variation in radii with a single tool which minimizes the need for a large tool bit inventory and frequent tool bit changes and which is simple in design, versatile and is easy to use.\nIn view of the above mentioned problems and limitations associated with traditional cutting tools, it was recognized by the present inventor that there is an unfulfilled need for an improved cutting tool which can be varied in size to accommodate changing radii needs on a part being machined and one which is simple in design, practical, fun to use and is economically manufactured.\nAccordingly, it becomes clear that there is a great need for a variable radii cutting tool which overcomes the disadvantages associated with traditional tool bits used in machining operations. Such a cutting tool should be one that minimizes the need to frequently change tool bits and one which eliminates the fixed radii features associated with standard tool bits."}
{"text": "As users become more active on computers, their concerns about identity protection increase. Users are concerned about the security of their digital identity. There are numerous tools to protect various aspects of identity. However, in the prior art, each of these aspects of a user's digital identity must be individually handled. The user can download anti-virus software, encryption software, and other tools to attempt to protect their system and identity.\nOne of the biggest concerns for users is passwords used to access accounts ranging from 401 K accounts, to bank accounts, and email. Stories about users' passwords being stolen or compromised are abundant. Users want privacy of their account, but they also want convenience. With the huge numbers of accounts that most users have these days, there is often a problem with choosing passwords for each of them, and ensuring that those passwords are of sufficient complexity and changed regularly. This causes many millions of dollars of loss, to consumers as well as banks and other institutions."}
{"text": "The present invention relates generally to testing integrated circuits, and more particularly to built-in-test circuits and output response analyzers for integrated circuits.\nIntegrated circuits are typically tested multiple times while they are manufactured. Often, individual circuits are tested while they are part of a wafer, which contain thousands of integrated circuits. Nonfunctional die are identified, for example with an ink spot, during a test referred to as wafer sort. After wafer sort, the die are separated and packaged. The packaged devices are testing again—this is referred to as final test. Additional testing may be done, for example sample devices may be tested under extreme environmental conditions.\nDuring these tests, test data, also referred to as test vectors, which typically include data and clock signals, are provided to the integrated circuit by a tester. The input test data may be generated by a circuit or software test pattern generator. Conventionally the integrated circuit operates on the input test data and provides output test data back to the tester. An output response analyzer in the tester checks the output test data for errors, and passes or rejects the device.\nIt is desirable to test each node in an integrated circuit. However, integrated circuits are becoming extremely complicated and may include hundreds of thousands of logic elements. At the same time, it is desirable to reduce the number of pins on the device in order to simplify device packaging and reduce printed circuit board complexity and space. The result is that many internal nodes on integrated circuits are difficult to reach electrically by device pins.\nAccordingly, it is desirable to include test circuitry on the integrated circuit itself, such that these internal nodes may be more thoroughly tested. Further, it is desirable to provide an internal test circuit that is capable of testing using test patterns other than simple all ones or all zeros patterns. Also, it is desirable to be able to perform such tests without the addition of complicated circuitry. It is also desirable that the internal circuitry require no or a limited number of pins, such that device pin count may be maintained."}
{"text": "1. Technical Field\nThe present invention relates to a circuit board, an electro-optical device, and an electric apparatus.\n2. Related Art\nRecently, as a device substituting for a thin-film field effect transistor using inorganic materials such as silicon, an organic thin-film field effect transistor has attracted attention. This is because a method of manufacturing constituent elements of such an organic thin-film field effect transistor is simple, properties of the constituent material can be changed by changing the molecular structure of the organic semiconductor material and the organic thin film field effect transistor is more flexible, lighter, and stronger than the thin-film field effect transistor formed using an inorganic semiconductor.\nSince there is interaction between a substrate and a semiconductor layer forming a transistor, it is important which material is selected for forming the substrate. For example, JP-A-2005-101555 discloses a study of an organic field effect transistor having an organic insulating layer containing an aromatic compound.\nHowever, the optimum amount of an active group to be contained in an insulating layer of a transistor and a relationship of electrical resistance between the insulating layer and a semiconductor layer required consideration. Additionally, mobility of charge carriers of the semiconductor may not be sufficiently improved."}
{"text": "1. Field of Invention\nHop flavors, compositions thereof, unreduced alpha and isoalpha acids, reduced alpha and isoalpha acids, removal of undesirable odor-forming impurities therefrom, hop flavors which do not contain such impurities, flavoring of beer including ale therewith.\n2. Prior Art and the Invention Broadly Stated\nHop flavors of the isoalpha acid types enumerated under FIELD OF INVENTION have attracted widespread interest in recent years for use in the controlled and standardized flavoring of beer including ale, especially by post-fermentation treatment. The present invention involves the discovery that certain hitherto-unrecognized undesirable odor-forming impurities (the exact structure and nature of which is yet unknown) exist in all of these hop flavors which, if allowed to remain, in time give rise to undesirable aromas of a fruity, estery, fatty or sour nature upon standinq or storage of beers flavored post-fermentation therewith. The present invention accordingly also involves the removal of such undesirable odor-forming impurities from the starting hop flavor with production of a purified hop flavor which is essentially free of undesirable amounts of such odor-forming impurities and hence also free of such odor- or aroma-contaminating tendency when employed to flavor beer post-fermentation. The method for eliminating the said impurities involves their extraction from the hop flavor into water at a pH above about 4, preferably above about 6 or 7, the exact preference in pH depending upon the acidity of the hop flavor being extracted, sometimes in the presence of solvent such as residual carbon dioxide, hexane, or lower-alkanol such as ethanol, and preferably in the presence of a salt for purposes of keeping most if not all of the hop flavor itself in the organic phase, an important economic consideration. Although the PRIOR ART has employed highly-acidic (pH 2 and below) aqueous extractions of such starting hop flavors to remove previously-recognized impurities such as metals, tannins, and sugars, extractions at such low pH levels will not and does not remove the odor-forming impurities discovered and eliminated according to the present invention. Neither does the procedure of the prior art, which has been used to extract crude isomerized whole hop extracts with water to produce an aqueous phase, still containing the impurities together with small amounts of isoacids which will dissolve in the water at the pH level employed, eliminate these impurities or provide a suitable post-fermentation flavor, even though such solutions have been suggested for flavoring of beer post-fermentation and are reported to be non-haze-forming. (Westermann U.S. Pat. No. 3,798,332)."}
{"text": "1. Field of the Invention\nThe present invention relates to a reflective collection-type light receiving unit and a light receiving apparatus for spatial light communications using the reflective collection-type light receiving unit and more particularly to a reflective collection-type light receiving unit for spatial light communications for performing communications using visible light or infrared light radiated in a space and a light receiving apparatus for spatial light communications in which a plurality of reflective collection-type light receiving units are arranged in an array to realize communications higher in sensitivity and responsiveness.\n2. Description of Related Art\nWireless communications by using radio waves as a communication medium have found many uses in various fields such as cellular phone networks, wireless LANs and short-distance wireless communications.\nHowever, in wireless communications in which radio waves are used as a medium, it is impossible to increase the transmission power, with the influence of electromagnetic waves on the human body taken into account, where the transmitting and receiving are performed near people. Further, since the frequency range of radio waves used in wireless communications has already been allocated in many areas and actually used, it is impossible to use freely frequencies over a wide range. Still further, under special circumstances such as medical institutions, there are some restrictions, for example, the limited use of radio waves.\nThus, in recent years, there has been developed visible light communications in which visible light is used as a communication medium, and a visible light communication system has been proposed in Japanese Published Patent Application No. 2007-266794 A.\nIn this type of a conventional visible light communication system, ordinarily, a light receiving element equipped with an optical system including a collecting lens is used as a light receiving unit on the side of a receiver, and light projected from a projector on the side of a transmitter is collected by the optical system and focused on the light receiving element, thereby outputting a light receiving signal. In particular, in visible light communications movable in a relative position between the transmitter and the receiver, in order to increase a light collecting percentage of light, with consideration given to a case where light radiated from the projector of the transmitter is lower in the amount of light arriving at a light receiving unit of the receiver, the optical system including the collecting lens is required to be larger in size, the light receiving unit is made relatively larger in configuration, and in particular the use of the collecting lens increases the thickness dimension of the light receiving unit, thus posing a problem that it is difficult to downsize the optical system.\nOn the other hand, Japanese Published Patent Application No. H10-322276 A has disclosed a light receiving unit for spatial light communications in which a reflector lens is arranged on a light receiving path of a light receiving element to decrease the thickness dimension of the light receiving unit. However, in this reflective collection-type light receiving unit, a convex mirror and a concave mirror are arranged inside the reflector lens to constitute an optical system and a light receiving element is arranged behind the reflector lens, thereby posing another problem that the optical system is increased in size.\nFurther, the light receiving unit of this type of a conventional spatial light communication apparatus is constituted in such a manner that a plurality of reflective collection-type light receiving units are arranged in an array and a light receiving element of each of the reflective collection-type light receiving units is ordinarily mounted on a substrate of a light receiving circuit. In particular, where the plurality of light receiving elements are connected in parallel to the input side of an amplifying part of the light receiving circuit, a floating electrostatic capacitance generated between a pair of electrode leads of each light receiving element and the ground is multiplied only by the number of light receiving elements to be used and increased accordingly. There has been posed such a problem that an electrostatic capacitance generated on the electrode leads of the anode and the cathode of the light receiving element is increased to reduce a response speed of the light receiving element, thereby inhibiting spatial light communications at high speed."}
{"text": "There exist a variety of portable combustion devices that may be used, for example, as a stove for cooking or heating, among other uses. In the past, portable combustion devices required a variety of fuels such as those used for liquid fuel stoves, portable and fixed wood stoves and compressed gas fuel stoves. These stoves were used in a variety of different situations such as for camping, emergency or rescue situations, during a power outage or a similar scenario when traditional larger-scale cooking sources (for example gas/electric cooking rages and barbeque grilles) are not available.\nWhile previous devices did provide a combustion device, they all had the similar drawback that they all required a special fuel source. Whether the fuel was kerosene, gasoline, propane or a similar fuel source they all required some canister of fuel to be purchased and carried along with the combustion device. These canisters may be a single use canisters or multi-use canister. However, once the fuel is expended, an additional fuel canister must be supplied to maintain combustion. The availability of fuel poses a serious drawback should the canister run out, especially if used for hiking, or emergency situations. Likely, it is impossible to obtain a new canister of gas or liquid fuel while out camping or during an emergency, and even if it is possible, there is no guarantee that the person would be able to find replacements for their type of fuel and canister. Therefore the operator is required to bring extra canisters of fuel, which adds more weight, and the canisters become yet another accessory that must be carried in addition to device. Moreover, should a person need to use their portable combustion device for an extended length of time, carrying a multitude of canisters for fuel becomes impractical.\nAdditionally, the canisters of fuel must be disposed of properly after use. Since the canisters are often under pressure and most fuels are harmful or toxic, the canisters can not be left behind at a campsite, or during an emergency scenario. In some cases the canister may not be thrown away like other types of trash, but need to be disposed in a special manner. With a variety of fuels, canister types, and portable combustion devices, which are all made by different companies a person would be limited in their options and be forced to continue to use a single fuel source.\nAside from the traditional campfire, there have been several attempts to overcome the disadvantages associated with fuel-burning cooking/heating appliances, one such attempt being a biomass combustion device or biomass stove. A biomass stove is able to burn a variety of biomass fuels such that a person would not be required to supply a compressed liquid fuel source for their portable combustion device. The user could employ whatever fuel is available to achieve the same combustion results as the previous devices. Additionally, the user need not dispose of potentially toxic canisters, and he or she would not be required to carry the appliance's fuel source with him or her. Also, whatever fuel is unused can simply be left behind.\nIn fact, in many developing countries some form of biomass combustion or biomass stove, burning wood chips, twigs, leaves, peat, etc. (or low-grade coal in some instances) is often used as the primary cooking device for the family meal. Generally the biomass combustion is carried out in some form of simple fire-driven stove, which is used for both heating and cooking. However, biomass devices typically have drawbacks that render them undesirable in many situations. For example, biomass stoves are not “clean” burning, that is they produce soot and smoke during the combustion process, which can cover pots, pans, or even food during the combustion. Also, the smoke generated during the combustion process can make the biomass stoves potentially dangerous for indoor use—as the smoke contains large quantities of soot and other toxic combustion products, including carbon monoxide. Nevertheless, such stoves are often used in small indoor spaces with inadequate ventilation in developing countries. In addition, many biomass stoves are heavy or built-in, making them less suitable for field use, as they are not easily portable.\nOne technique for increasing the efficiency of combustion is to provide a driven airflow through the biomass with a powered fan, using a variety of motive power sources including, but not limited to, spring and clockwork mechanisms, compressed air/gas and electricity. In the field, the fan is typically driven by a battery or other electrical source. However, batteries may have a short life in use, and/or be discharged (dead), or be in need of recharging, when needed. In addition, batteries are expensive and often unavailable to peoples of developing countries. In addition, batteries are environmentally unfriendly and often disposed of improperly.\nAccordingly, it is desirable to provide a portable combustion device that is capable of being used as a stove for cooking and/or heating, which burns commonly available wood and other biomass, and provides similar heat output as other liquid or gas fuels, without the need for batteries or disposable/refillable canisters of fuel. The stove can desirably be easy to carry, low-maintenance and burn fuel without exhausting significant soot, smoke or toxic combustion byproducts. The stove can desirably employ inexpensive and commercially available components and conventional construction techniques in its manufacture."}
{"text": "The present invention relates to video amplifiers and, more particularly, to an active impedance-matched video amplifier."}
{"text": "Robotic welding assemblies are commonly used to weld manufactured components, such as vehicle components. The weld tips used in these assemblies must be periodically dressed to remove excess weld material that tends to buildup with repeated usage. Removal of the weld material is necessary to maintain the precise weld tip geometry necessary for accurately positioned welds of reliable strength. Weld tip dressers are commonly used to remove the weld material. A weld tip dresser has a cutting tool that acts on the weld tips to remove the excess material and reconfigure the tip to the desired geometry. Typically, the material removed by the weld tip dresser falls to the factory floor and must be periodically swept away or otherwise disposed of. Minimizing the downtime required for the weld tip dressing procedure and associated cleanup is desirable to maintain scheduled production goals."}
{"text": "Various physical parameters are monitored and/or measured during operation of an aircraft. Various methods of monitoring and/or measuring these physical parameters can be used. For example, sensors and/or transducers can measure position and/or relative locations of members of an aircraft. Some sensors and/or transducers can be affixed to an aircraft at specific locations so as to produce signals indicative of various physical phenomena experienced at those specific locations. These signals can then be transmitted to an analyzer that receives and interprets the signals. Some specific locations where it would be desirable to affix a sensor and/or transducer might be locations that have harsh environments. For example, some such locations might expose an affixed sensor to high temperatures, high pressures, high levels of exposure to electromagnetic interference, etc.\nIn many of these harsh environment locations, optical transducers have found use. Optical sensors and/or transducers can produce optical signals indicative of various physical phenomena. For example, optical sensors and/or transducers can produce optical signals indicative of stress, strain, temperature, tilt, rotation, vibration, pressure, etc. Various sensors and/or transducers employ various types of technologies. For example, some sensors use Fabry-Pérot Interferometry (FPI). Some sensors use Fiber Bragg Grating (FBG) technologies. Some sensors use intensity modulation techniques. Some sensors use various optical diffraction gratings to generate an output signal. Some of these technologies and techniques produce optical signals having a narrow-band spectrum or a specific wavelength that is indicative of the measured parameter. Spectrum analysis and/or spectral measurement of such signals can be performed to determine a measure of the physical phenomena indicated by the narrow-band spectrum or specific wavelength of the optical signal.\nPhotodetectors used in the determination of wavelength detection required precise temperature characterization. Such precise temperature characterization both increased the complexity of wavelength measurement systems and ultimately limited the fidelity of such systems. The fidelity of the spectral measurement can correlate to the precision and/or accuracy of the measured physical parameter. The frequency rate at which such spectral measurements can be performed can limit the frequency that measurements of such physical parameters can be provided. Thus, the present disclosure is directed to methods and systems for fast determination of the narrow-band spectrum or specific wavelength of such optical signals, thereby obtaining fast determination of the physical parameters measured by the various optical sensors."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of the read error processing in magnetic disk drives. More specifically, the invention relates to a method of preventing read errors resulting from the thermal asperity (TA) of the servo area of the disk.\n2. Background Information\nA magnetic disk drive is a drive which writes data to the disk surface or reads data from the disk surface by a change of magnetization. A converter is positioned over a desired position on the track on which data has been recorded, and data is read out of the track of the disk rotating at high speed, or data is written to the track. The magnetic head of the magnetic disk drive is slightly spaced from the disk surface and positioned over a predetermined track. Magneto-resistive (MR) head has been adopted as one of the (magnetic field) conversion heads. The output resistance of the MR head changes as the magnetic field changes. This resistance change with a predetermined current to pass through the MR element converts MR resistance change to voltage signal, as reading out data.\nHowever, the method of reading out this resistance change has a thermal asperity problem as one of its operational problems. The thermal asperity can be defined as follows. The large asperity of the disk surface touches a magnetic head and crashes, and temperature change causes resistance change to occur in the MR element. Because of this resistance change, an abnormal signal will occur.\nThe phenomenon of a thermal asperity often causes a strip temperature to rise to more than 100 degrees locally. The cause of this temperature rise is a mechanical crash between the head portion, including the MR stripe, and the asperity of the disk surface. Since the (normal) resistance change of the magnetic head by the magnetization change in the medium (disk), which is caused by normal reading, is less than 1% of the MR stripe resistance, the signal change resulting from the temperature rise, which is caused by the thermal asperity, largely exceeds the resistance change caused by a normal read signal, so normal data reading is hindered.\nIn conventional methods which cope with the thermal asperity with respect to data that is read out, when a sharp signal change, which would be determined to be a thermal asperity, is read out, an appropriate alternation is added to that read signal so that it can be processed as data, or that read signal is corrected with an error correction code. However, these conventional methods have the disadvantage that external additional hardware is required, or in the case of a large burst error which should be with an error correction code, correction becomes impossible and finally it should be processed as an unreadable hard error. Thus, the conventional methods cannot solve the thermal asperity problem sufficiently."}
{"text": "Many passenger lifts designed to carry people are primarily used by skiers and are in use for a short period of the year only. Since many very charming mountain resorts have their tourist season only during the winter months and since the local lifts and inns are only in full use during this short period of the year, they tend not to be especially profitable.\nPlaygrounds have common children's slides which are used almost all year long. However, such slides are mounted on a structure which limits the total length of the slide.\nThe present invention is designed to increase the economy of passenger lifts, e.g., ski lifts. This may be done by building a long mountain slide at a ski slope, which slide consists of several straight and curved segments joined together and descending the entire mountain. The possibility of a unique downhill ride on the new slide will, of course, attract people to use the \"ski\" lift during the summer or non-snowy seasons.\nIn accordance with the invention, a simple and suitable support for the mountain slide, and an appropriate mat to sit on while sliding down an extended distance, as well as device to return the mat to the top of the mountain all have been provided as part of the new recreational facility.\nThe new slide is erected at the slope and may be as the particular incline permits, e.g., a slide may be 1,000 meters long. The slide may be in operation during the entire year, and, therefore, it adds to the year-round operation of any lift and its economy.\nThe partly straight and partly curved segments of the new slide may be arranged in serpentines, whereby the segments with narrow curves have higher edges. An appropriate material for construction of the slide segments is asbestos concrete. This material is inexpensive, durable in bad weather, shock-resistant, and it has good gliding characteristics after treatment with wax or a similar material. Preferably, each segment has at its upper part an overlap reaching under the next segment going uphill. In case of settling of one segment or its footing, no step will be formed against which a passenger might slide, although a small drop may be created.\nThe sliding surface of the segments is hollowed out or dished and is bordered lengthwise by open bottom flanges arching upwards and with their bottom edges extending outwardly away from the sliding surface. This design makes the segments so solid that no longitudinal supports are required. Furthermore, the hollowed out sliding surface permits a pleasant and comfortable ride down the new slide.\nThe footings for the new slide normally consist of ties which rest on bars or stakes driven into or otherwise anchored to the ground. If the ground conditions permit, it is also possible to attach the ties to posts. Alternatively, the posts may be driven into the ground for use as footings. This is recommended when the slide travels over a depression in the terrain. When erecting a slide or part of it on a rocky surface, the footings may be in the form of concrete foundations.\nAccording to another characteristic of the invention, the slide segments are attached to the footings by adjustable supports in order to compensate for any difference in height of the footings. Preferably, the segments rest with their arched edges on adjustable struts having rounded upper portions. Short pieces of pipe with a diameter approximately equal to that of the radius of the arched slide segment edges are suited for the support of the arched edges of the slide segments. The bearing surface of the slide will then be sufficiently supported to avoid any breaking of the asbestos concrete slabs. The struts are arranged on the footings in such a fashion that they swivel in all directions in order to fit any position, i.e., they are \"universal\", thus compensating for deviations in nominal height and lateral positioning of the footings.\nIn order to avoid undue movements in a completed slide and to secure slide segments against lifting off their anchorages, the strut supports may be specially connected with the flanges of the slide segments in accordance with another characteristic of the invention. This special support may be achieved by directing each strut through an opening in a butt strap fastened onto the slide flange, the apertured leg of the butt strap being positioned directly beneath the rounded supporting piece. The aperture diameter of the strap is larger than the strut diameter in order to allow for a limited movement of the slide in relation to the footing. With the butt strap arranged underneath the footing, the slide cannot lift off the strut, and, consequently, it cannot lift off the footing.\nIn accordance with a more specific detail of the invention, each strut is in the form of a threaded rod and the lower end (opposite the rounded upper supporting piece) thereof is threaded into a base tube provided with a mating thread. By selective rotation of the base tube, the total length of the strut may be changed to compensate for inaccuracies, deviations in height, etc. The base tube extends through an opening into a yoke-like fixture attached to the footing, and a lock inside the yoke prevents lifting off, e.g., a retaining ring may be arranged in a ring groove in the base tube for this purpose. The base tube has, at the end opposite the thread (its lower end), a convex bearing surface which always nests solidly on the footing if the strut orientation should vary slightly.\nIn accordance with the invention, in order to prevent the adjustable struts from improperly tilting sideways or downhill, the upper part of each strut may be connected by a stay or bracing with an adjacent footing or with the other end of the footing on which it rests. The stay or bracing is adjustable for precise alignment and orientation of each strut.\nAccording to another characteristic of the invention, there are new and improved sleds for sliding down the new slide. These sleds may consist simply of mats with handles on each side. The handles may be simply leather straps suitably attached to the mats with tubular rivets, screws, or the like. As will be understood, the mats protect the clothing of passengers and, furthermore, prevent the marring of the slides which may be caused by shoes, particularly by heels. For the securing of feet and the eliminating of soiling of the slide, the mats may have foothold devices, such as pockets attached to the mats. By attaching ropes to the mats near the foothold devices, a passenger who has fallen backwards may readily straighten up by using the ropes.\nFurthermore, the mats may be equipped with a brake shoe located at a central rear portion. In order to slow the sled down, the brakes may be pressed down. Specifically, the mats may have an indentation to which the brake shoe is attached. If a passenger leans back while sliding, it is inevitable that the weight of the body will push down on the brake shoe, thus establishing braking contact with the slide and reducing the speed of the sled automatically.\nFor comparatively short slides, the mat, as described previously will suffice. These mats are not especially suitable for long slides with curves, where downhill sled speed has to be substantially reduced in the curves in order to meet safety requirements. Furthermore, the conditions of long slides are more likely to be affected by the weather, e.g., sections exposed to the sun have different friction properties than those in the shade. Recommended for such cases are sleds having their brake shoes attached to one end of a swiveling lever, the other end of which forms a handle. The brake may be activated at will to effect the amount of deceleration required.\nThe sled may comprise a frame provided on the bottom with glide linings, which frame forms a seat and accommodates the brake lever. The brake lever is mounted along the longitudinal axis of the sled and, therefore, in longitudinal direction of the frame with its handle end extending beyond the frame. The brake lever is pivoted on the frame between its handle end and brake shoe end.\nThe rear end (in the downhill or sliding direction) of the brake handle has a plate accommodating an interchangeable brake pad, preferably consisting of rubber. In order to improve braking properties, the bottom of the brake pad includes indentations comparable to the treads of car tires. The front (in sliding direction) of the brake plate is generally pointed, is downwardly bent, and is downwardly biased by a spring supported by the brake handle. The bent tip of the brake plate engages the slide surface, regardless of the angle of inclination of the brake handle due to the spring pressure, and deflects dirt and moisture, so that the effectiveness of the brake shoe is not reduced.\nThe top of the sled frame is equipped with a seat covered by a mat of felt or similar material, preferably extending over the frame in both longitudinal directions of the sled. The bottom of the frame is equipped with glide linings. These may consist of interchangeable plastic plates with plastic needles or bristles at the bottom. Plastic skids may also be utilized to advantage.\nFor a more complete understanding of the present invention and a better appreciation of its attendant advantages, reference should be made to the following description taken in conjunction with the accompanying drawings."}
{"text": "At present, digital-data recordable media, such as digital audio tapes (DATs), compact discs (CDs), and digital video discs (DVDs), are distributed, and various types of content such as music data and picture data are recorded as digital data in the media and are widely distributed.\nSuch digital data differs from analog data, and is free from data deterioration due to data copying between recording media. If copying is limitlessly permitted, there is a possibility that the rights of a content-copyright holder and other content-related-right holders may be violated. For protecting the right of the digital data, there is the SCMS (Serial Copy Management System) as a copyright protection technology.\nThe SCMS (Serial Copy Management System) is a digital-data-copying restriction system. It allows only-one-generation (1-Generation) copying, and prohibits digital copying for two or more generations. Specifically, by recording, in digital-data-recorded media, a code representing copying only once, restriction of copying is performed based on the code.\nNevertheless, in one-generation copying control using the SCMS, based on the recorded code in the media which represents copying only once, that is, the bit state, it is determined whether or not copying is allowed. Accordingly, by using a device that can freely operate the bit, rewriting of the code is made possible, and many copies identical to the original data can be made. Therefore, in particular, PC-used (personal computer-used) copying of digital data such as CDs, which is free from restrictions of law, is actually free.\nIn addition, for systems for the purpose of protecting copyright that records/plays back content such as pictures and music, a system has been proposed in which content is encrypted and provided to a user and in which a key for decryption is provided to a normal user.\nBy way of example, there is a system configuration in which various types of content, such as music data, picture data, and game programs which are encrypted, are distributed to users by using the Internet or media such as CDs and DVDs and in which only a person identified as a normal user is provided with a means for decrypting the encrypted content, that is, a decryption key.\nThe encrypted data can be returned to usable decrypted data (plaintext) by decryption processing based on a predetermined procedure. Such a data encryption/decryption method is conventionally known in which an encryption key is used for information-encrypting processing and a decryption key is used for decrypting processing.\nAmong various types of examples of data encryption/decryption methods using an encryption key and a decryption key, there is a method as an example that is a so-called a common key cryptosystem. In the common key cryptosystem, by setting an encryption key for data-encrypting processing and a decryption key for data decryption to be common, and providing a normal user with a common key for the encryption processing and decryption, data accessing by a user having no key is excluded. A typical of this system is the DES (Data encryption standard)\nThe encryption key and the decryption key for the above encryption processing and decryption can be obtained by using a unidirectional function, such as the Hash function, based on, for example, a password or the like. The unidirectional function is a function in which reverse finding of its input from its output is very difficult. For example, by using a user-decided password as an input in an application of the unidirectional function, an encryption key and a decryption key are generated based on the output. It is substantially impossible to perform reverse finding of the password as the original data from the encryption key and the decryption key obtained as described above.\nA system in which a process using the encryption key for use in encryption and a process using the decryption key for use in decryption have different algorithms is a so-called public key cryptosystem. The public key cryptosystem is a system in which unspecified users use an usable public key, and an encrypted document for a specified person is encryption-processed by using a public key issued by the specified person. The document encrypted by the public key becomes able to be decryption-processed by using only a secret key corresponding to the public key used in the encryption process. Since a secret key is possessed by a person who issues a public key, a document encrypted by the public key can be decrypted by only the person who possesses the secret key. One typical public key cryptosystem is the RSA (Rivest-Shamir-Adelman) cryptography. Use of such a cryptosystem enables a system in which encrypted content can be decrypted only for a normal user.\nIn this system, for example, a 2-bit EMI (Encryption Mode Indicator) is defined as copy control information. When the EMI is 00B (B indicates that the value before it is a binary number), it indicates that content is of a Copy-freely type, and when the EMI is 01B, it indicates that content is of a No-more-copies type in which the content may further not be copied. When the EMI is 10B, it indicates that content is of a Copy-one-generation type in which copying only once is allowed, and when the EMI is 11B, it indicates that content is a Copy-never type in which copying is prohibited.\nWhen the EMI represents the Copy-freely or Copy-one-generation type, it is determined that content can be copied. Alternatively, when the EMI represents the No-more-copies or Copy-never type, it is determined that content cannot be copied. If management of the copy rule information is appropriately executed, copyright protection is realized.\nHowever, even in the content providing system using encryption, if information on copying rules which is recorded on a medium such as a CD or a DVD is rewritten by an invalid user, a problem occurs in that copying ignoring the original copying rules becomes executable."}
{"text": "Heating of engine and vehicle components increases engine efficiency by reducing the viscosity of various lubricating fluids, thus reducing the energy needed to pump the fluids among the vehicle components. Further, a heated engine converts fuel to shaft work more efficiently than does a cold engine because of lowered heat losses to the combustion chamber. While the engine may be heated relatively rapidly via the heat produced during combustion, transferring the heat to downstream driveline components may take a longer duration, thus delaying when the driveline components reach operating temperature and extending the period that the vehicle operates with reduced efficiency. Further, heating the driveline components via energy sourced from fuel or other stored potential energy reduces fuel economy.\nThe inventors have recognized the above issues and provide a method to at least partly address them. In one embodiment, a method for a vehicle comprises heating a fluid with kinetic vehicle energy in response to a vehicle braking request, and directing the fluid to a driveline component.\nThe kinetic vehicle energy produced in response to a vehicle braking request may be used to heat a fluid, and this heated fluid may in turn heat a driveline component. To heat a fluid via kinetic vehicle energy, a viscous brake may be engaged responsive to the vehicle braking request. The viscous brake may decelerate the vehicle as requested by the vehicle braking request, while simultaneously heating the fluid used to lubricate the viscous brake. The heated fluid is then used to heat the driveline components, via a heat exchanger for example. In this way, the driveline components may be heated without relying solely on combustion heat, improving engine efficiency and fuel economy.\nThe above advantages and other advantages, and features of the present description will be readily apparent from the following Detailed Description when taken alone or in connection with the accompanying drawings.\nIt should be understood that the summary above is provided to introduce in simplified form a selection of concepts that are further described in the detailed description. It is not meant to identify key or essential features of the claimed subject matter, the scope of which is defined uniquely by the claims that follow the detailed description. Furthermore, the claimed subject matter is not limited to implementations that solve any disadvantages noted above or in any part of this disclosure."}
{"text": "Low grade heat, which is typically considered less than 100 degrees, represents a significant waste energy stream in industrial processes, power generation and transport applications. Recovery and re-use of such waste streams is desirable. An example of a technology which has been proposed for this purpose is a Thermoelectric Generator (TEG). Unfortunately, TEGs are relatively expensive. Another largely experimental approach that has been proposed to recover such energy employs Shape-Memory Alloys.\nA Shape-Memory Alloy (SMA) is an alloy that “remembers” its original, cold-forged shape which, once deformed, returns to its pre-deformed shape upon heating. This material is a lightweight, solid-state alternative to conventional actuators such as hydraulic, pneumatic, and motor-based systems.\nThe three main types of Shape-Memory Alloys are the copper-zinc-aluminium-nickel, copper-aluminium-nickel, and nickel-titanium (NiTi) alloys but SMAs can also be created, for example, by alloying zinc, copper, gold and iron. The list is non-exhaustive.\nThe memory of such materials has been employed or proposed since the early 1970s for use in heat recovery processes and in particular by constructing SMA engines which recover energy from heat as motion. Recent publications relating to energy recovery devices include PCT Patent Publication number WO2013/087490, assigned to the assignee of the present invention. Other patent publications include US2005/150223 (United Technologies) and U.S. 2013/341845 (Zanella). It is desirable to translate the contraction of the SMA or NTE material into a mechanical force in an efficient manner. It is not a trivial task and generally is complicated and involves significant energy losses.\nIt is therefore an object to provide an improved system and method in an energy recovery device."}
{"text": "1. Field of the Invention\nThe present invention relates to a lithographic apparatus, a projection assembly and a combination of a structure and an active damping assembly.\n2. Background Art\nA lithographic apparatus is a machine that applies a desired pattern onto a substrate, usually onto a target portion of the substrate. A lithographic apparatus can be used, for example, in the manufacture of integrated circuits (ICs). In such a case, a patterning device, which is alternatively referred to as a mask or a reticle, may be used to generate a circuit pattern to be formed on an individual layer of the IC. This pattern can be transferred onto a target portion (e.g. including part of, one, or several dies) on a substrate (e.g. a silicon wafer). Transfer of the pattern is typically via imaging onto a layer of radiation-sensitive material (resist) provided on the substrate. In general, a single substrate will contain a network of adjacent target portions that are successively patterned. Conventional lithographic apparatus include so-called steppers, in which each target portion is irradiated by exposing an entire pattern onto the target portion at once, and so-called scanners, in which each target portion is irradiated by scanning the pattern through a radiation beam in a given direction (the “scanning”-direction) while synchronously scanning the substrate parallel or anti-parallel to this direction. It is also possible to transfer the pattern from the patterning device to the substrate by imprinting the pattern onto the substrate.\nThe high accuracy and high resolution currently aimed at in lithography require an accurate positioning of parts of the lithographic apparatus such as the reticle stage to hold the mask, the projection system and the substrate table to hold the substrate, with respect to each other. Apart from the positioning of e.g. the reticle stage and the substrate table, this also poses requirements on the projection system. The projection system in current implementations may consist of a carrying structure, such as a lens mount (in case of transmissive optics) or a mirror frame (in case of reflective optics) and a plurality of optical elements such as lens elements, mirrors, etc. In operation, the projection system may be subject to vibrations due to a plurality of causes. As an example, movements of parts in the lithographic apparatus may result in vibrations of a frame to which the projection system is attached, a movement of a stage such as the substrate stage or the reticle stage, or accelerations/decelerations thereof, which may result in a gas stream and/or turbulence and/or acoustic waves affecting the projection system. Such disturbances may result in vibrations of the projection system as a whole or of parts thereof. By such vibrations, displacements of lens elements or mirrors may be caused, which may in turn result in an imaging error, i.e. an error in the projection of the pattern on the substrate.\nCommonly, a damping system is provided to dampen vibrations of the projection system or parts thereof. Thereto, a passive damping system may be provided as known in many forms, or a active damping system, or a combination of a passive and a active damping system. In this document, the term active damping system is to be understood as a damping system which comprises a sensor to detect an effect of a vibration (e.g. a position sensor, velocity sensor, acceleration sensor, etc) and a actuator to act on the structure to be damped or a part thereof, the actuator being driven by e.g. a controller in dependency of a signal provided by the sensor. By driving the actuator in dependency of the actuator, i.e. in response to a signal provided by the sensor, an effect of the vibration on the projection system or a part thereof, may be reduced or cancelled to a certain extent. An example of such active damping system may be provided by a feedback loop: the sensor to provide a position quantity (such as a position, speed, acceleration, jerk, etc of the projection system or a part thereof), the controller being provided with the position quantity and generating a controller output signal to drive the actuator, the actuator in turn acting on the projecting system or the part thereof so that a feedback loop is provided. The controller may be formed by any type of controller and may be implemented in the software to be executed by a microprocessor, microcontroller, or any other programmable device, or may be implemented by dedicated hardware.\nA problem needing to be solved relates to the stability of the feedback loop, i.e. it is desirable to achieve a frequency behavior of the feedback loop wherein ringing and/or oscillation is prevented. At the same time, a high bandwidth of the active damping system is required, as a high bandwidth of the active damping system will allow to suppress vibrations within such high bandwidth. Due to the ever increasing demands on speed of the lithographic apparatus, movements in the lithographic apparatus tend to take place at a higher speed and consequently involving faster transients, which may result in a generation of vibrations at increasingly higher frequencies. Thus, as speed increases, it is desirable to achieve a higher bandwidth of the active damping system.\nA projection system is commonly built up from a variety of parts, comprising e.g. lenses, mirrors and/or other optical elements, lens mountings and/or mirror mountings, a housing of the projection system such as a lens body, etc. As a consequence, a frequency behavior of the projection system starts, at a low frequency extreme, as a rigid body mass, thereby providing a transfer function which is inversely proportional to frequency, as illustrated in FIG. 2, where on a vertical axis a transfer function has been plotted (on a logarithmic scale), while in a horizontal axis a frequency has been plotted (on a logarithmic scale), and where RBM indicates a frequency behavior substantially corresponding to a rigid body mass. In a resonance frequency range RES, a resonance of the projection system is observed, which may be followed by a plurality of further resonance frequencies with increasing frequency, thereby overall resulting in an increase of the magnitude of the transfer function. Effectively, as from the resonance area, the projection system does not behave as a single object anymore, however instead shows a variety of resonance phenomena each corresponding to resonance of an element of the projection system. As a result thereof, the higher the frequency, the lower the remaining mass which “contributes” to the transfer function, which may be considered an explanation for the fact that the magnitude of the transfer function increases with increasing frequency, in the frequency range above the resonance frequency range. As will be understood by a skilled person, the frequency behavior of the projection system as outlined above, may result in stability problems when attempting to achieve a bandwidth of the active damping system which reaches or exceeds the resonance of the projection system. The transfer function may be expressed in terms of e.g. velocity of the projection system as a function of a force on the projection system. It is noted that the transfer function may also be expressed in any other suitable quantity, such as acceleration of the projection system as a result of force on projection system. In that case, a low frequency behavior of the transfer function will show to be frequency independent, followed by a resonance frequency range and an increase of the transfer function (showing multiple resonance peaks) above the resonance frequency range."}
{"text": "1. Field of the invention\nThe present invention relates to a solar battery type indication apparatus used in road signs, guide signs, or the like, and particularly relates to a solar battery type indication apparatus which is improved to attain reduction both in weight and in cost and eliminate the unevenness of bending of optical fibers to thereby make it possible to perform light transmission securely, receive various parts in a body of the indication apparatus efficiently and connect the optical fibers easily without entanglement in the indication apparatus body.\nFurther, the present invention relates to an indication apparatus for performing an indication by using light from light-emitting diodes (LEDs), a light source such as a solar lighting apparatus, or the like, and particularly relates to an indication apparatus which is improved in optical fiber-mounting property.\n2. Description of the Related Art"}
{"text": "1. Field of the Invention\nThe present invention relates to electrochemical electrodes and, more particularly, to combination electrodes for measuring the ion concentration of solutions.\n2. Description of the Prior Art\nCombination electrodes comprising a sensing portion and a reference portion incorporated in a glass or a plastic structure are known in the art. One known electrode assembly for measuring pH includes a generally flat, pH-sensitive glass membrane having one face for contacting a sample to be measured and an opposite face contacting an internal electrolyte within the assembly. In addition, a liquid junction is provided between the sample and a reference electrolyte within the assembly. By virtue of its flat ion-sensing structure the electrode assembly is able to measure a small quantity of sample solution on a relatively flat surface without the need to immerse the electrode in a large volume of the solution. Thus, the electrode is particularly suited for measuring, for example, the pH of moisture on the skin of a human subject.\nWhile electrodes of the foregoing nature have performed satisfactorily, they have been confined generally to laboratory, clinical, or other specialized applications requiring trained operating personnel. Further, the sensing portion of such prior electrode assemblies must be carefully positioned on the sample to be measured with the electrode assembly vertically oriented with respect to the sample. This is necessary to ensure that the conventional liquid junction included in the prior art electrode makes contact with the sample and to ensure that any air bubbles within the electrode assembly rise to the top thereof so as not to contact either the ion-sensitive membrane or the liquid junction. If these precautions are not taken, erroneous measurements can result. In addition, the reference electrolyte is generally consumed during normal operation of the electrode and must be periodically refilled. The refilling operation requires a skilled operator to ensure that a proper electrolyte is employed and that no contamination is introduced during the filling operation.\nRecently, the need has arisen for a combination electrode assembly which can be operated over prolonged time periods at any physical orientation by non-technical operators. For example, a particular need exists in the cosmetic and hairdressing industry for an electrode assembly which can be routinely employed by hairdressers, beauty operators, and the like for monitoring the pH of a subject's hair, skin, or other body parts. The present invention provides such an assembly."}
{"text": "A semiconductor gas sensor is operated to determine a concentration of a detection gas by using a degree by which a resistance of a detection material such as a metal oxide, a polymer and a carbon nano tube is changed during an oxidation/reduction reaction of the detection gas and the detection material. To this end, the semiconductor gas sensor employs a dedicated heater for a smooth oxidation/reduction reaction of a detection gas and a detection material, and such a reaction is generated in a temperature section of 200 to 500 degrees Celsius.\nAn electrochemical gas sensor is operated to determine a concentration of a gas through a degree by which a current or an electromotive force is changed while the detection gas and a detection material such as an alkaline metal or an alkaline earth metal carbonate react each other at a high temperature to influence an ion conductivity of a solid electrolyte.\nThe semiconductor gas sensor or the electrochemical gas sensor employs a heater therein in terms of operation principle. The detection material used for the semiconductor gas sensor or the electrochemical gas sensor mainly includes a metal oxide, an alkaline metal carbonate, a carbon nano tube, graphene or a polymer.\nHowever, a gas sensor according to the related art is known that sensitivity, selectivity and the like are influenced by humidity, that is, water molecules in the air (Jianwei Gong et al., Sensors and Actuators B 114(2006) 32-39). To solve this problem, a study for improving humidity stability of a gas sensor through a filter for adsorbing moisture to a gas inlet port or removing moisture from the gas inlet port while the gas sensor is packaged (U.S. Pat. No. 5,841,021). Such a filter technology helps remove moisture of a part of a gas sensor and secure stability of the gas sensor, deteriorates efficiency of the gas sensor in a situation where excessive moisture is present. And the filter has a limit in life when the filter is stored for a long time. Therefore, a process of replacing a filter after a predetermined time period is required to protect a gas sensor."}
{"text": "1. The Field of the Disclosure\nThe disclosure relates generally to a system and methods for determining whether an energetic substance or material has experienced a reaction (“go”) or a non-reaction (“no-go”) for storage, transportation or in-process handling, and more particularly, but not necessarily entirely, to a system and methods using a video capturing device, a CPU or computer, sensitivity test equipment, such as an electrostatic discharge device or impact assessment device for testing and assessing the substance or material reaction or explosion sensitivities, and a set of rules or instructions to be followed for quantifying and determining whether a reaction has occurred or not.\n2. Description of Related Art\nReaction detection for sensitivity test equipment is not automated. As a consequence, the determination of whether a reaction actually occurred (i.e., a go reaction) or whether a no reaction occurred (i.e., a no-go reaction) is subjective in nature and is based on an operator's varying experience and perception of what actually transpired during a reaction event. Because of the subjectivity of determining whether a reaction occurred or not, standardization of sensitivity test results between laboratories is very difficult to achieve. Standardization between laboratories is possible provided an objective system and methods for determining whether a reaction occurred or not are used.\nIn the industry, various types of equipment are used to assist an operator in determining whether a reaction occurred or not while testing an energetic material or substance for sensitivity. Some of these operator assisted devices include, but are not limited to, a noise dosimeter, a gas analyzer, a light meter, a strain gauge, and a video capturing device (such as a standard or high-speed camera). One of the best ways in the industry to determine whether a reaction occurred or not is to use a high-speed video capture device for at least the following reasons: (1) it provides a reviewable visual record of the reaction event; and (2) it provides better spatial and temporal visual resolution of the reaction event.\nThe disclosure is directed to a unique and advantageous system and methods for determining whether a reaction occurred or not using automated equipment, thereby reducing the amount of subjectivity resulting from an operator determination based on an image or collection of images. The automated system and method may use a high-speed video capturing device, sensitivity test equipment, a computer processor and a set of rules or instructions that objectively compares a set of quantified image trial data to a set of quantified image baseline or background data to determine whether a reaction occurred.\nAn additional difference between the disclosure and what is done in the industry is instead of the operator making a subjective reaction determination based on sensory perception or a less subjective determination using a threshold for auditory, or a gas analyzer, or intensity of light, the operator makes a determination on quantifiable data relating to the acceptable level of error or the likelihood of a false positive and/or a false negative. With that level of error identified, the reaction detection threshold can then be calculated through the use of the system and method of the disclosure. Knowledge of the error or the likelihood of a false positive and/or a false negative allows for much better risk assessment of the sensitivity testing outcome.\nThe disclosure improves upon known techniques used in the industry by, inter alia, using four unique identifiers or quantifiers relating to the images identified as being of interest and significant. Those identifiers or quantifiers include: brightness, shape, buoyancy, and the uniformity of the event. The disclosure also detects decomposition or reaction of energetic materials. The method and process disclosed may also be advantageous in that multiple characteristics of the images are simultaneously quantified to determine if a reaction has occurred. It is noteworthy that none of the known systems or methods known to Applicant provides the above-identified advantages.\nThe features and advantages of the disclosure will be set forth in the description which follows, and in part will be apparent from the description, or may be learned by the practice of the disclosure without undue experimentation. The features and advantages of the disclosure may be realized and obtained by the methods and means of the instruments and combinations particularly pointed out in the appended claims. Any discussion of documents, acts, materials, devices, articles or the like which has been included in the present specification is not to be taken as an admission that any or all of these matters form part of the prior art base, or were common general knowledge in the field relevant to the disclosure as it existed before the priority date of each claim of this application."}
{"text": "1. Technical Field\nThis invention relates generally to textile sleeves, and more particularly to braided textile sleeves.\n2. Related Art\nIt is known to protect elongate members in textile sleeves against a variety of environmental conditions and affects, or to just contain elongate members in textile sleeves for bundling and routing purposes, such as in knit, woven or braided sleeves. In the case of braided sleeves, the braided wall is commonly braided as a circumferentially continuous, seamless wall, sometimes referred to as a ‘closed’ wall. One known advantage of a closed, braided wall construction is that the wall can be circumferentially expanded to facilitate sliding the wall over an elongated member by manually pushing and physically holding the opposite ends of the wall in a compressed fashion. By pushing the opposite ends toward one another and manually holding the wall in an axially compressed state, the braided wall is caused to take on an increased diameter and a reduced length. When in the increased diameter state, the wall can be readily disposed over the elongate member. Then, after sleeve is installed over the elongate member, the installer can release the wall and the opposite ends automatically spring axially away from one another, thereby taking on a circumferentially decreased diameter and increased length.\nAlthough the aforementioned ability to increase and decrease the diameter of a braided wall has an advantage over some other known types of sleeve construction, such as woven sleeves, it does come with potential drawbacks. Namely, the ability to manually increase the diameter of the braided sleeve requires applying a continual, externally applied compression force during installation, which can prove challenging, and thus, can complicate the ability of the installer to readily install the sleeve over the elongate member. Further complicating installation of a braided sleeve arises when the sleeve has a relatively long length. With the sleeve having a relatively long length, difficultly arises by having to axially compress the opposite ends toward one another without causing the sleeve to fold or buckle along the length of the sleeve. In addition, upon releasing the wall to have the sleeve resume its lengthened, decreased diameter state, the wall generally has a tendency to spring back, at least partially, toward its axially compressed configuration due to a pattern retention phenomenon caused by friction between the interlaced yarns. As such, the effective length of the sleeve can be unintentionally decreased.\nIt is further known to stretch, cut or tear openings in a sidewall of the sleeve to provide access for branches of wires and branched wire connections, such as with Y or T-shaped connectors. Unfortunately, stretching yarns to form openings is time consuming and can result in inadvertently breaking yarns, while cutting the sidewall typically results in frayed, loose yarns, which not only appear unsightly, but can also diminish the useful life of the sleeve by allowing the opening to become enlarged over time, thereby subjecting the contents to contamination. It is further known to bundle branches and branched connections with multiple sleeves, thereby avoiding having to form openings in a sleeve wall, and further, to use tape about the branches. Although this can prove useful, it is time consuming, cumbersome, and as discussed above, can result in an unappealing appearance, particularly if the tape comes loose. In addition, increased cost is generally associated with having to use multiple sleeves and/or tape."}
{"text": "1. Field of the Invention\nThe present invention relates to laser scanning systems for reading bar code symbols, and more particularly, pertains to a novel digitizer processing circuit for detecting the times at which a scanning light spot moves from a space onto a bar code symbol or moves from a bar code signal onto a space.\n2. Description of the Prior Art\nThe increased use of bar code symbols to identify products, particularly in retail businesses, has resulted in the development of various bar code reading systems. Many users of bar code readers require portable hand-held scanners which place a premium on small size, lightweight, and low power consumption requirements for the devices. One such system is a laser scanning bar code reading system as described in U.S. Pat. No. 4,496,831, commonly assigned to the same assignee as the present patent application.\nThe laser scanning system disclosed in U.S. Pat. No. 4,496,831 includes a portable hand-held scanning head which may be embodied in various shapes but preferably has a gun-shaped housing made of lightweight plastic. A handle and barrel portion are provided to house the various components of the scanning head therein. Within the barrel portion are mounted a miniature light source, a miniature optic train including focusing lenses and a scanning system for directing light from the light source across a bar code symbol, and miniature sensing means for detecting reflected light from the bar code symbol being scanned.\nThe miniature light source can comprise a laser tube such as a coaxial helium neon laser tube, or preferably a semiconductor laser diode which is considerably smaller and lighter than a laser tube, thus reducing the required size and weight of the scanning head and making the scanning head easier to handle and more maneuverable. Light generated by the light source passes through the optic train which directs the beam onto the scanning system which is mounted within the barrel portion of the scanning head. The scanning system sweeps the laser beam across the bar code symbol and comprises at least one scanning motor for sweeping the beam lengthwise across the symbol, and may comprise two motors wherein the second motor sweeps the beam widthwise across the symbol. Light reflecting means such as mirrors are mounted on the motor shafts to direct the beam through the outlet port to the symbol. A sensing circuit then detects and processes the light reflected from the symbol and generally comprises a photodetector element such as a semiconductor photodiode. The user positions the hand-held unit so the scan pattern traverses the symbol to be read, and the photodetector element produces serial electrical signals to be processed for identifying the bar code. A circuit such as that of the present invention for a digitizer signal processing circuit for a bar code produces a signal which is directed to a bar pattern decoder circuit for decoding the bar pattern. The reader unit can have a movable trigger employed to allow the user to activate the light beam and detector circuitry when pointed at the symbol to be read, thereby conserving battery life if the unit is self-powered. The lightweight plastic housing contains the laser light source, the detector, the optics, signal processing circuitry, a CPU, and a battery. The reader is designed to be aimed at a bar code symbol by the user from a position where the reader is spaced from the symbol, i.e., not touching the symbol or moving across the symbol. Typically, this type of hand-held bar code reader is specified to operate in the range of perhaps several inches.\nAlthough the present invention is described with respect to linear or single line bar codes, it is not limited to such embodiments, and may also be applicable to more complex scanning patterns and to stacked or two dimensional bar codes such as Code 49 and similar symbologies. The present invention may also find application for use with various machine vision or optical character recognition applications in which information is derived from other types of indicia such as characters or from the surface characteristics of the article being scanned.\nIt is well known that it is desirable to locate in time the moment when the center of the moving scanning laser spot is located on the edge of a bar as it moves onto the bar and again when it moves off the bar. It is also well known that there is an inflection point in the analog signal that corresponds to the above-stated moments in time. Various forms of signal processing have been used in bar code scanners to locate these inflection points.\nThe most common technical approach to accomplish this is to use a zero crossing detector to locate when the second derivative of the analog signal crosses through zero, as this zero crossing corresponds in time to the inflection points of the analog signal. One problem with this approach is that there also are zero crossings of the second derivative that do not correspond to the inflection point of the analog signal. Various technical approaches have, therefore, been developed to distinguish between the desired zero crossings and false ones. Most of these approaches depend upon information obtained from either the analog signal or the first derivative of the analog signal. Examples of these technical approaches are disclosed in Spectra Physics Patent Nos. 4,000,397 and 4,749,879. Another known method is to offset the analog signal by half of its amplitude so that it is centered around zero and then to compare it directly to the second derivative with a voltage comparator. Another common feature is a circuit to ensure noise-free margins the white area on either side of the bar code) by setting a noise threshold based upon the averaging of noise over a period of time. Signals below this average noise level will not be digitized."}
{"text": "1. Field of the Invention\nApplicant's invention relates to surgical instruments and, more particularly, to trocars. Trocars are used to pierce or puncture an anatomical cavity to provide communication with the inside of the cavity during a surgical procedure.\n2. Background Information\nEndoscopic surgery, particularly laparoscopic surgery, is currently becoming a significant method for performing surgeries. It is projected that by the year, 2000 half of all surgical procedures will be performed endoscopically. Laparoscopic surgery has become the surgical procedure of choice because of its patient care advantages over \"open surgery.\"\nFor the past several decades, endoscopic surgery has been available as a method of diagnosis and, for a very limited number of disorders, a treatment. Until recently, a factor limiting the types of surgeries that could be performed laparoscopically was the ability to employ intraoperative assistance. In the past, endoscopes allowed only direct visualization by the surgeon, such as the endoscope disclosed in U.S. Pat. No. 4,254,762 issued to Yoon. This led to the situation where the surgeon had one hand holding the laparoscope to his eye and then had only one hand available to operate.\nFortunately, miniaturization of video camera computer chips has led to the development of video cameras that can easily be attached to an endoscope or laparoscope. During surgery, connecting a video camera and monitor to the laparoscope enables all the operating room personnel to view the surgical procedure, rather than just the surgeon. Thus, the operating room personnel are able to provide operative assistance just as they do with open surgery. The type and number of surgical procedures amenable to laparoscopic surgery is presently one of the most rapidly developing areas of medicine.\nThe pivotal advantage of laparoscopic surgery over open surgery is the decreased post-operative recovery time. In many instances, a patient is able to leave the hospital within twenty four hours after laparoscopic surgery has been performed. This is compared to a five day to ten day hospitalization necessary to recover from an open surgical procedure. Additionally, laparoscopic surgery provides a decreased incidence of post-operative abdominal adhesions and decreased post-operative pain with enhanced cosmetic results.\nAn essential medical instrument for endoscopic procedures is the trocar. Trocars are sharp, pointed surgical instruments used to puncture the wall of an anatomical cavity. The trocar consists of a tube or cannula and a cutting element called an obturator or stylet. The obturator fits within the cannula and has a sharp piercing tip at its end.\nA conventional laparoscopic trocar insertion procedure usually follows insufflation of the abdominal cavity with CO.sub.2 gas. The introduction of CO.sub.2 gas into the abdominal cavity lifts the abdominal wall away from the internal viscera. Once this is done, the abdominal wall is penetrated with the trocar. After insertion of the trocar through the abdominal wall, the surgeon removes the obturator leaving the cannula or tube protruding through the body wall. A laparoscope or laparoscopic instruments can then be inserted through the cannula to view internal organs or perform surgical procedures.\nPenetrating the wall of the abdominal cavity with the trocar is done quickly. The sharp point of the obturator encounters great resistance from the skin, muscle, and tissue membranes of the abdominal wall while it is being pushed through these structures. Once the trocar's sharp point and blade pass through the abdominal wall and into the cavity, the resistance to the trocar drops quickly. Unless the surgeon immediately stops pushing the trocar just as soon as penetration of the abdominal wall is complete, there is a chance that the trocar will penetrate further into the abdominal cavity and injure internal organs.\nWithin the abdominal cavity, the obturator's sharp point could easily injure or cut an internal organ upon the slightest contact. If an internal organ is inadvertently injured or cut, unless immediate and massive hemorrhage occurs, the injury may not become apparent until long after completion of the surgery. At a minimum, such an injury will delay a patient's recovery and, more likely, could seriously endanger the patient's health. Additional corrective surgery on an open basis may be required, subjecting the patient to additional risks and costs.\nPrior to 1987, the only trocars available for laparoscopic use were instruments made from stainless steel, such as those disclosed in U.S. Pat. No. 3,994,287 issued to Turp et al., and U.S. Pat. No 3,613,684 issued to Sheridan. A problem common to all of these \"classic\" trocars is that they do not have a safety shield which covers the sharp, cutting tip of the obturator once it pierces the cavity wall.\nSeveral changes and additions have been made on the functional design of these classic trocars. The most significant improvement on the classic trocar is the addition of a spring-loaded safety shield that snaps forward to cover the sharp point and blade of the obturator once the trocar has penetrated the abdominal wall, such as those disclosed in U.S. Pat. No. 4,601,710 issued to Moll, U.S. Pat. No. 4,654,030 issued to Moll et al. and U.S. Pat. No. 4,535,773 issued to Yoon. In these devices, the safety shield is a plastic sleeve which is concentrically mounted to the obturator. Because of this safety feature, trocars with a spring loaded safety shield have become the most used trocars in laparoscopic surgery.\nHowever, these spring loaded safety shields have cumbersome safety shield control mechanisms which are difficult to tell if the safety shield is armed or engaged. To be sure of the safety shield's operation, a surgeon will need to verify the proper operation of the safety shield prior to use. Since the safety shields are mounted to the obturator, surgeons are required to test the safety shield's operation with the obturator in place by manually pressing the safety shield of the trocar. Unfortunately, this results in many slight puncture wounds being experienced by the surgeons as they are attempting to verify the safety shield's operation.\nCurrently, laparoscopic trocars with the spring loaded safety shields are manufactured only as a thin wall plastic disposable instruments. These light weight plastic instruments are used once and then discarded as medical waste, adding significantly to the already escalating health care costs. For example, the single use disposable plastic trocars cost approximately sixty five dollars to ninety dollars each. Usually two to four trocars are used for each laparoscopic procedure. Thus, surgical costs are unnecessarily increased about two hundred dollars to three hundred and fifty dollars per laparoscopic procedure, as well as adding to the overbearing problem of medical waste by the use of these disposable trocars. Presently, however, the increased health care cost has been unavoidable because the only safety shielded trocar available was a disposable item. No one had invented a easily disassembled, easily cleaned, sterilized, and easily reassembled for reuse, safety shielded trocar."}
{"text": "An image recognition and classification system (a machine vision system) includes a preprocessor in which a \"top-down\" method is used to extract features from an image, an associative learning neural network system which groups the features into patterns and classifies the patterns, and an attentional mechanism which focuses additional preprocessing and a neural network on relevant parts of an image.\nAttempts to recognize and classify images have led to construction of automated artificial machine vision systems and to development of strategies to learn patterns in images and to recognize and classify images by using the learned patterns. Those developing artificial systems have continually attempted to incorporate principles of biological systems into their strategies, because biological systems outperform all artificial systems, implemented or proposed, by a wide margin. For example, machine vision systems based on artificial neural networks have been implemented on digital parallel computers, but a parallel implementation only provides an increase in speed without an increase in performance. Thus, the goal of emulating the pattern recognition performance of biological systems still eludes computer scientists.\nIn order for a biological nervous system to discriminate objects two fundamental problems must be solved: object segmentation and binding. \"Object segmentation\" deals with distinguishing separate objects; \"binding\" deals with how specific attributes such as shape and depth, are linked to create an individual object. A question addressed by object segmentation mechanisms is to which overlapping object does a border belong? An image of an object may be occluded (divided) by an overlapping image, and will need to be reconstructed as a whole image. Models have been proposed to explain this process, for example, using artificial neural networks. (Sajda and Finkel, 1992)\nDuring the past half century, the theoretical infrastructure of machine vision systems has developed both top-down (beginning with large features of the image) and bottom-up (beginning at the lowest level of resolution, usually a pixel) views. However, actual development has focused almost exclusively on bottom-up approaches as exemplified by the title of Pentland's illuminating book From Pixels to Predicates, and comments therein such as: \"Processing is primarily data-driven (i.e., bottom-up), although it can be responsive to the goals and expectations at the higher levels.\" (Pentland, 1986, part 1, page 1).\nOngoing efforts have focused on the extraction of \"features\" in an image by local manipulations of small micro-features (often 3.times.3 rarely more than 9.times.9 pixel areas), with the intent of identifying larger features (macro-features) from their combination. The paucity of robust results from this approach may be attributed to several causes, two of the most important of which are (1) that the mathematical operations performed on the small areas are usually differential operators such as edge detectors that enhance rather than reduce noise; and (2) that not even humans are very good at visual recognition when allowed only a small instantaneous field of view. Similarity of tactual and visual picture recognition with limited field of view. Loomis et al. (1991).\nDuring this same time period, cognitive psychologists and neurobiologists have made impressive advances in research on the processing mechanisms that are at work in the visual cortex of mammals, particularly cats and monkeys. Electrophysiological and psychophysical experiments on cats and monkeys demonstrate a wide variety of feature selective cells in the visual cortex. In the mammalian cortex, these include simple cells (Hubel and Weisel, 1962), whose shape is closely approximated by a Gabor function (Daugman, 1985; Jones and Palmer, 1987) or a difference of Gaussian functions; end-stopped cells (often called first order hypercomplex cells) (Hubel & Weisel, 1965; Gilbert, 1977); color sensitive cells; and even cells that respond only to faces. (Desimone, 1991). Face-selective cells in the temporal cortex of monkeys. Desimone (1991).\nComplex cells and second order hypercomplex cells (Hubel and Weisel, 1962, 1965) are sensitive to the same features as simple and first order hypercomplex cells, respectively. One of the differences among these cells, of interest in the context of feature extraction from static images, is that the complex and second order hypercomplex cells have larger receptive fields than simple cells, and are insensitive to location of micro-features within their receptive fields.\nIn the development of artificial systems, preprocessing of data derived from an image has been used to extract features from an image and to select features for further processing by machine vision systems. Preprocessing generally proceeds in steps from the \"bottom-up,\" although \"top-down\" preprocessing has been suggested as a model for human vision. Preprocessing is accomplished by preprocessors, which may be implemented in hardware or software. In some systems, preprocessors have served as the first layer of a two layered neural network. Preprocessing strategies have included subdividing a whole image to be processed into sub-images. Various filters have been suggested to operate on the data, converting the data to a different form or value distribution. Control masks have been used to focus a network on a specific domain of an image.\nPrior approaches to the problem of modeling biological preprocessing have been addressed by Grossberg (1988) and Fukushima (1988). For example, the neocognitron neural network developed by Fukushima conceptually models simple, complex, first order and second order hypercomplex cells as well as layers of cells that are sensitive to higher order features. Second order hypercomplex cells are constructed from combinations of first order hypercomplex cells; complex cells are constructed from combinations of simple cells, and the like.\nOne means of making a complex cell insensitive to location, the approach used by Fukushima, is to design it to receive input from several adjacent simple cells, whose frequency and orientation tuning are similar. The complex cell is made sensitive enough to respond when only one of the simple cells responds to a stimulus. The result is a complex cell with the same frequency and orientation tuning as the simple cells, whose receptive field size is equivalent to the total receptive field size of all its input simple cells combined. Furthermore, the complex cell is insensitive to where in its receptive field the luminance pattern is located (i.e., the complex cell is insensitive to which simple cell has been activated.) Trying to apply biological principles to artificial vision systems, Porat and Zeevi (1989) determined from their work and the work of others, that \"primitives of image representations in vision have a wavelet form similar to Gabor elementary functions (EF's),\" and proposed a method for texture discrimination in images using a Gabor approach.\nAlthough Porat and Zeevi (1989) proposed that \"These localized operators (referring to Gabor functions) are also suitable for a pyramidal scheme of multiresolution which appears to be characteristic of vision, and can also serve as oriented-edge operators and in pattern recognition tasks,\" (p. 116), they adopted the prevailing approach to the process as a bottom-up hierarchy.\nAn alternative to extracting features using predefined, generally applicable fixed filters (detectors), such as generated by Gabor and end-stop filters, is to design a system that generates its own feature detectors. In biological systems, the feature detectors must be general enough to handle all possible inputs encountered during the life experiences of the animal. It has been shown that a linear neural network with a correlation rule, when stimulated by random noise, will develop feature detectors similar to the center-surround and Gabor filters found in some artificial visual systems. However, in most practical applications of artificial networks, the universe of possible inputs is more restricted. This suggests that a system for adaptive filter generation that can develop feature detectors specific to the range of images that are encountered in a practical application would be highly desirable. Self-modifying learning algorithms have been pursued wherein a learning algorithm learns about its own effectiveness and modifies itself so that it is the most effective algorithm for solving a certain class of problems.\nDespite extensive efforts and much progress, \"Forty years of research in artificial neural networks has yielded networks with the neural complexity of, perhaps, a sea slug.\" (Wenskay, 1991) Image recognition and classification remains a major frontier. The present invention advances toward this frontier."}
{"text": "The endogenous cholinergic neurotransmitter, acetylcholine, exert its biological effect via two types of cholinergic receptors, the muscarinic Acetyl Choline Receptors (mAChR) and the nicotinic Acetyl Choline Receptors (nAChR).\nAs it is well established that muscarinic acetylcholine receptors dominate quantitatively over nicotinic acetylcholine receptors in the brain area important to memory and cognition, and much research aimed at the development of agents for the treatment of memory related disorders have focused on the synthesis of muscarinic acetylcholine receptor modulators.\nRecently, however, an interest in the development of nAChR modulators has emerged. Several diseases are associated with degeneration of the cholinergic system i.e. senile dementia of the Alzheimer type, vascular dementia and cognitive impairment due to the organic brain damage disease related directly to alcoholism. Indeed several CNS disorders can be attributed to a cholinergic deficiency, a dopaminergic deficiency, an adrenergic deficiency or a serotonergic deficiency.\nWO 2004/029053 describes 1,4-diazabicycloalkane derivatives useful as modulators of the nicotinic receptors. However, the 1,4-diaza-bicyclo[3.2.2]nonyl oxadiazolyl derivatives of the present invention have not been described."}
{"text": "The present invention relates to a novel concept in genetic engineering, according to which a eukaryotic organism is transformed and crossed so as to provide an offspring including a first exogene in a first chromosome of a chromosome pair and a second exogene in a second chromosome of the chromosome pair, the exogenes exhibiting allelic relationship, such that the exogenes obligatorily segregate to different gametes. The present invention further relates to expression cassettes for implementing the novel concept, to methods of implementing same for the obtainment of reversible male sterility in plants, and further to plants and plant products obtained thereby.\nThe gene is the basic functional unit of heredity. The gene concept was first set forth by Mendel in 1865 as a result of experiments conducted on pea plants. Mendel proposed that a genetic determinant of a specific character is passed on from one generation to the next as a unit without any blending of the units, a theory now known as the “one gene one trait” theory.\nTwo of mendel's proposed laws serve as the basis for modern genetics. Mendel's first law states that two members (alleles) of a gene pair obligatorily segregate to different gametes, such that one half of the gametes carry one member (allele) of the pair and the other half of the gametes carry the other member (allele) of the gene pair. Mendel's second law states that during gamete formation the segregation of alleles of a gene pair is independent of alleles other gene pairs.\nFollowing Mendel's research, the discovery of chromosomes and the existence of linkage between genes closely positioned on a single chromosome have further contributed to our knowledge of inherited traits. Nowadays it is known that the second law of Mendel is valid for genes residing on different chromosomes and for genes residing on the same chromosome, provided that they are at least 50 centiMorgan apart, whereas for closer genes residing on the same chromosome, linkage dis-equilibrium is experienced.\nThus, in an individual, a genetic trait in it's simplest form is determined by a pair of closely related sequences known as alleles. Each allele, which is optionally an alternative form of a gene, resides on one of the chromosomes of a chromosome pair. The pair of alleles of a given individual can be identical in sequence or they can differ in sequence to various degree. In meiosis, which is the process of cell division which leads to the formation of sex cells (gametes), the chromosomes of any given chromosome pair segregate into different gametes. This chromosomal segregation is superseded by rearrangement of the DNA in chromosome pairs which takes place via inter chromosomal recombination events.\nA specific trait of an individual is determined by a specific combination of two alleles carried by the individual. Thus, traits of an individual are determined by the genetic material inherited by that individual from both parents as determined by chromosomal segregation and recombination events in each of the individual's parents.\nAlleles carried on a chromosome pair encode the same genetic trait. Allelic variance ensures a high degree of distinction between individuals.\nTransformation of an exogene to a genome is typically random. Assuming one site of exogene integration, a resultant genome is said to be heterozygous of the type 1/0. By careful breeding and selection, a homozygous state for the exogene, i.e., 1/1 is obtainable. Although one can transform a genome with a second exogene (2), it is practically impossible, using conventional transformation techniques, to direct the second exogene to a given integration site, in a different genome. As such, using such conventional molecular biology and breeding techniques, it is practically impossible to achieve exogene allelism of the type 1/2.\nAlthough it is theoretically possible to direct an exogene to a specific chromosomal site by gene knock-in techniques, which therefore could have been used to generate exogene allelism of the type 1/2, such techniques, are very poor in transformation yields. Indeed, no attempts of generating exogene allelism were made so far.\nThe production of hybrid plants has been practiced ever since the beginning of this century. Plant breeders discovered that crossing two distinct parental lines often resulted in a hybrid plant which displays what is termed as hybrid vigor or heterosis, which is characterized by an increased crop yield and/or adaptation to both biotic and abiotic stresses. Early on, hybrid plant production concentrated mainly on corn and sorghum which were found to benefit from crossing of well compatible parental lines. Later it was discovered that hybrid plant production is also applicable to other plant species.\nSeveral explanations have been proposed for the existence of heterosis. It is clear that a wider allelic variety exist in hybrids plant lines as is compared to inbred plant lines, since the alleles present in a hybrid plant are inherited from two distinct parent plant lines. This allelic variety is favored as a likely explanation for hybrid vigor, however, at present no scientific evidence has been brought forth to support this explanation.\nSince hybrid plants have been demonstrated to be superior to inbred lines with respect to yield and vigor, the development of hybrid seeds is one of the prime objectives of the seed industry. In addition, since hybrid plant varieties result from a unique combination, the possibility of duplicating or reusing the hybrid seeds is minimized, thus, providing breeders with an inherent commercial protection.\nThe production of hybrid seed on a large scale is challenging because many crops have both male and female reproductive organs (stamen and pistil) on the same plant, either within a single flower or in separate flowers. This arrangement results in a high level of self pollination and makes large scale directed crosses between parental lines to generate hybrid plants difficult to accomplish.\nTo guarantee that out-crossing will occur during the production of hybrid seed, breeders have either manually or mechanically removed stamens from one parental line. Although such manual emasculation is effective for some plants, such as wheat, it is labor intensive and impractical for plants with small bi-sexual flowers. To cross-pollinate such plants breeders have often resorted to using naturally occurring male sterile mutants in efforts to produce hybrid seeds. Although the use of naturally occurring male sterile mutants enables to cross pollinate plants with small bi-sexual flowers, the availability of such mutants and oftentimes the poor genetic makeup thereof severely limited wide spread use of this approach. In addition, using sterile mutants for out-crossing typically results in retention of sterility in a large fraction of the produced seeds. Such seeds when sown would then produce sterile plants which in the case of self pollinator crops such as wheat would lead to no crop yield for a fraction of the plants. To overcome the resultant low crop yield, a breeder is forced to select out the sterile seeds which can be a nearly impossible task.\nSince important crops such as rice, wheat are self-pollinating plants with small bi-sexual flowers, there was a need and several attempts were made to develop systems for pollination control to assist in the production of F1 hybrids.\nHowever, several conditions must exist in order to obtain economically feasible F1 hybrid plants. The male sterile line must be 100% sterile and yet female fertile, a natural pollen transfer from the male fertile line to the male sterile line must be facilitated, and in cases of grain or fruit crops, male fertility restoration (MFR) should be enabled in order to obtain crop yield in the F1 progeny.\nThere exist several mechanisms of male sterility in plants. One such mechanism is cytoplasmic male sterility (CMS). In general CMS has been used in corn, rice, sorghum and onion on a limited basis, see, for example, U.S. Pat. No. 3,861,079 to Patterson.\nHowever, reliance on a single cytoplasmic-male-sterile system for the production of all hybrid plants is undesirable because it leaves the entire hybrid stock vulnerable to plant pathogens. For example, extensive use of one corn cytotype, cmsT lead to an eptiphytic outbreak of Southern Corn Leaf Blight in the early 1970's. Thus, it is important to develop alternative methods to produce male sterile lines in plant species where only a single male-sterility system is available.\nAn emerging alternative to cytoplasmic sterility is a nuclear sterility. Nuclear male-sterile-based pollination systems rely upon the introduction of a male sterility trait to one parental plant followed by the introduction of a fertility-restorer gene, as a result of cross-pollination, to produce fertile hybrid plants.\nTypically, genetically engineered nuclear male sterility is effected by expressing, in a controlled and targeted manner, a protein toxin which destroys anther tissues of the plant in which it is expressed. The expression of this toxin can be either silenced or antagonized in the F1 progeny by the introduction of a fertility-restorer gene, as a result of cross-pollination.\nMale reproductive processes in flowering plants occur in the anther. This organ is composed of several tissues and cell types, and is responsible for producing pollen grains that contain the sperm cells. A specialized anther tissue, the tapetum, plays an important role in pollen formation. The tapetum surrounds the pollen sac early in anther development, degenerates during the later stages of development and is not present as an organized tissue in the mature anther. The tapetum produces a number of proteins and other substances that either aid in pollen development or become components of the pollen outer wall. It is known that many male sterility mutations interfere with tapetum cell differentiation and/or function. Thus tapetal tissue is believed to be essential for the production of functional pollen grains.\nAs such, many nuclear male-sterile-based pollination systems known in the art utilize anther specific promoters to target the destruction of the tapetum or other anther specific tissues during maturation such that subsequent pollen development is arrested.\nFor example, U.S. Pat. Nos. 5,409,823, 5,659,124 and 5,824,542 to Crossland describe a dual method for producing male-sterile plants. Two genetically transformed plants, parents 1 and 2, are crossed to obtain male-sterile offspring. Parent 1 is transformed with an expression cassette comprising a nucleotide sequence encoding an anther-specific promoter which is operably linked to a nucleotide sequence encoding a trans-activator. Parent 2 is transformed with an expression cassette comprising a target nucleotide sequence, which is capable of being activated by the trans-activator, operably linked to a nucleotide sequence which encodes a toxin which can be an RNA or a polypeptide and which disrupts the formation of viable pollen. Therefore, crossing parent 1 with parent 2 results in male-sterile offspring. The male-sterile plants produced are useful for producing hybrid seed. Subsequent crossing of the male sterile plant with a plant which does not include the trans-activator restores fertility.\nAlthough this system produces desirable male sterility in a portion of the produced offspring, restoration of fertility which relies upon the segregation of the trans-activator gene from the toxin gene cannot be efficiently effected using this system. In the male sterile plant the segregation of these two genes during meiosis relies a great deal on their physical chromosomal location, as such when a male sterile plant is crossed with a plant which does not express the trans-activator, the trans-activator and the toxin can be provided by the male sterile parent as a result of which a portion of the produced progeny will retain undesirable male sterility.\nThus the method by Crossland is not efficiently applicable for self pollinators such as, for example, wheat, cotton and rice.\nU.S. Pat. No. 5,929,307 to Hodges describes a recombinant expression vector comprising a suicide gene flanked by site-specific recombination sequences, which vector, when introduced into a plant leads to male sterility. This invention further relates to a second expression vector including a recombinase gene which when introduced into the male sterile plant via a cross-pollinating restorer plant, leads to the deletion of the toxin gene and as a result to the restoration of fertility.\nIt will be appreciated in this case that since restoration of fertility is effected by a gene product which is introduced via a cross-pollinating restorer plant, a high copy number of the restorer gene must be maintained in the restorer plant such that highly efficient restoration of fertility in the resultant progeny is effected.\nThus, the presently available nuclear male-sterility pollination systems, although generally efficient in the production of male sterile plants, suffer from limitations when restoration of such male sterility via out-crossing is desirable, such as the case for grain or fruit producing hybrid plants. Since these systems cannot enable 100% restoration of fertility to the male sterile plants produced and since the resultant male sterile hybrid seeds are not typically selected out, a suboptimal crop yield from grain or fruit producing hybrid plants produced by these systems results.\nThere is thus a widely recognized need for, and it would be highly advantageous to have, a male sterile plant, which plant when out-crossed with a compatible male fertile plant generates male fertile offspring in substantially 100% of it's progeny. As further detailed hereinunder, such a plant can be generated by a transformation method which results in formation of allelism between two exogenes. There is thus a widely recognized need for, and it would be highly advantageous to have, a method and expression cassettes for the generation of exogene allelism."}
{"text": "1. Field of the Invention\nThe present invention generally relates to dispensing a product; and more particularly to a system and method of dispensing prepaid items using a uniquely identified container.\n2. Description of Background\nCurrently, a typical method for dispensing food items, such as soft drinks, is to sell a generic container to a customer who can then fill it at a freestanding kiosk. This method does not insure that the customer will not allow others to use the authorized container for refills, or that only the authorized container is used. This method also does not insure that the container cannot be used beyond the period of time for which the sale was intended.\nFor systems that dispense items other than soft drinks, for example use of a company gas card at a company gas pump, existing schemes do not insure that the receiving container (usually a vehicle) is authorized. Such is the case for when a company gas card is utilized to fill the gas tank of a personal automobile."}
{"text": "1. Technical Field\nThe present disclosure relates to sutures and more particularly to sutures made from, incorporating and/or coated with compositions of therapeutic absorbable glass. The present disclosure also relates to pledgets coated with compositions of therapeutic absorbable glass.\n2. Description of Related Art\nAntimicrobial agents have been associated with surgical devices to prevent contamination with germs. For example, U.S. Pat. No. 5,019,096 to Fox, Jr. et al. describes applying a coating to a medical device, the coating containing a matrix polymer and antimicrobial agents. Other examples of antimicrobial devices include U.S. Pat. Nos. 3,674,901; 3,705,938; 3,987,797; 4,024,871; and 4,612,337.\nSutures prepared from biocompatible bioabsorbable polymers are well known in the art and are described e.g., in U.S. Pat. Nos. 2,668,162; 2,703,316; 2,758,987; 3,225,766; 3,297,033; 3,422,181; 3,531,561; 3,565,077; 3,565,869; 3,620,218; 3,626,948; 3,636,956; 3,736,646; 3,772,420; 3,773,919; 3,792,010; 3,797,499; 3,839,297; 3,867,190; 3,878,284; 3,982,543; 4,047,533; 4,060,089; 4,137,921; 4,157,437; 4,234,775; 4,237,920; 4,300,565; 4,523,591, U.K. Patent No. 779,291; Gilding et al., Biocompatibility of Clinical Implant Materials, Vol. II, ch. 9: “Biodegradable Polymers” (1981). Synthetic biocompatible bioabsorbable multifilament sutures such as DEXON®, VICRYL®, and POLYSORB® commercially available from Ethicon, Inc. (Somerville, N.J.) and United States Surgical (Norwalk, Conn.) are well known in the industry.\nExamples of biocompatible non-bioabsorbable polymers useful for fabricating sutures include, but are not limited to, nylon, silk, polyester, polypropylene, polyethylene, cotton, linen, etc. Commercially available sutures fabricated from biocompatible non-bioabsorbable polymers, e.g., a polyester suture (SURGIDAC®, United States Surgical, Norwalk, Conn.) and a polyester braided suture (TICRON®, David & Geck, Danbury, Conn.) are also well known in the industry.\nSuture coating compositions are also well known in the art. For example, U.S. Pat. No. 4,027,676 describes an absorbable coating composition for sutures. Other suture coatings are described, e.g., in U.S. Pat. Nos. 4,624,256; 4,190,720; 4,582,052; 4,605,730; 4,700,704; 4,705,820; 4,788,979; 4,791,929; 4,994,074; 5,047,048; 5,100,433; 5,352,515; 5,032,638; 4,711,241; 4,705,820; and 4,201,216.\nWater-soluble glasses have been utilized for a variety of medical, cosmetic and other purposes. For example, UK Patent Specifications Nos. 1,565,906, 2,079,152, 2,077,585 and 2,146,531, describe the dissolution of glasses impregnated with various agents such as drugs, hormones, insecticides, spermicides, and fungicides to provide controlled release of these agents. The glass can be in the form of an implant or bolus. WO 98/44965, describes a water-soluble biodegradable glass composition containing various active agents, e.g., antimicrobials such as antibiotics and metal compounds, e.g., silver oxide, silver orthophosphate, steroids, painkillers, etc., which is used for implantation in soft tissue.\nWO 96/24364, describes a controlled release glass having various metals, e.g., silver, copper, and zinc useful for combating infections.\nWO 98/54104, describes the preparation of water-soluble glass fibres which optionally contain silver compounds. The glass fibres are utilized in undefined orthopaedic implants and tissue engineering applications.\nU.S. Pat. No. 5,470,585, describes the use of medicinal substances for topical applications, e.g., wound dressings, which include a water-soluble glass containing a silver compound.\nU.S. Pat. No. 4,612,923 is directed to glass-filled absorbable surgical devices made of a synthetic absorbable polymer containing an absorbable glass filler.\nU.S. Pat. Nos. 5,290,544 and 5,766,611, describe cosmetic products containing soluble glass which contains silver, copper or zinc ions having antibacterial activity. U.S. Pat. No. 5,330,770, describes a boron-free water-soluble glass containing silver oxide useful as a water treating agent.\nWhile the aforementioned references describe the use of water-soluble glass for certain implant and cosmetic applications, wound dressings, treating infections, and water-treating agents there is no indication in the references of water-soluble glass-based coating compositions for sutures or pledgets, or sutures manufactured from water-soluble glass, which provide controlled, sustained release of a therapeutic agent."}
{"text": "The invention relates to the field of computers, and more specifically to computer buses.\nPresent day computers comprise bus systems, onto which different devices may be plugged. More specifically, a bus system is often comprised of a bus controller and of a bus connected to the memory controller. Different devices may be connected to the bus, so as to be accessed by the bus controller.\nOne example of such buses is the DRAM bus designed by Rambus Inc. This bus is used for managing high speed DRAM devices. FIG. 1 is a schematic view of the architecture of this type of bus. It shows the memory controller 1, and the Rambus Channel 2. Several Direct Rambus DRAMs or Direct RDRams (trademark) 3-6 are connected to the Rambus Channel. As shown on FIG. 1, the memory controller as well as each of the RDRAMs comprises a Rambus interface 8 for using the bus. The bus 2 is terminated at one end by terminations, and is also connected to a reference voltage Vref as well as to a 400 MHz bus clock.\nAccording to the Rambus specification, there is also provided a power down mode; it is contemplated in the specification that the power down mode is used for reducing power consumption, notably in portable computers.\nRambus products sold on the market are organised in Rambus RIMM (trademark) memory modules, each module supporting 4, 6, 8, 12 or 16 Direct RDRAMs devices. RIMM modules are compatible with standard motherboard form factors; a motherboard usually supports up to three module sockets. The Direct Rambus Channel signals are daisy chained through each module. See Rambus RIMM Module Preliminary Information, document DL0078 available from Rambus Inc.\nThere is also provided in the Rambus specification a SPD (Serial Presence Detect) device. The purpose of the SPD is to store and provide sufficient information for a system to initialise the memory subsystem correctly: the SPD is a ROM device provided on each RIMM module, which includes information relating to the DRAM timing and device parameters, core organisation, module parameters, and other system level information. The SPD EEPROM devices of each RIMM module conform to the I2C wire protocol, and may be read into or written into by the memory controller of a Rambus system. See Direct Rambus SPD Specification 1.0, available from Rambus Inc.\nFIG. 2 is another view of a physical Rambus architecture, this time in an invalid configuration; it shows the memory controller 1 and the Direct Rambus Channel 2. Three modules 10-12 are connected to the bus; each side of each module may have up to 8 RDRAM devices, referenced again 3-6 on FIG. 2. Reference 13 is the SPD EEPROM of module 10; reference 14 shows the I2C protocol bus used by the memory controller for accessing the SPD EEPROMs of the different modules.\nMore details on Rambus may be found in the corresponding specification, issued by Rambus Inc. under the title Direct Rambus Technology Disclosure, Oct. 15, 1997.\nOne problem with Rambus is that the load of the bus is limited to 3 modules, and to 32 Direct RDRAM devices; if one of these limitations is violated, the bus system is not designed to be operational, and or even to boot at all; a computer in which the bus system is installed would in this case not be able to boot either.\nThis limitation on the number of modules is not likely in practice to be violated, since the bus normally comprises at most three module slots, and usually 2 or 3 module slots. However, a module may comprise up to 16 devices, so that the number of devices on the bus may exceed the highest allowable number of devices. This is the case in the configuration shown in FIG. 2.\nThus, the configuration of the bus hardware is such as to enable the bus to be improperly configured; in this case, a physical layer configuration constraint on the bus can be violated, and proper electrical operation of the bus is therefore not ensured. This can prevent the bus as a whole from booting properly. This possibility makes the system difficult for the user to upgrade or to diagnose problems that occur when they try to.\nA variety of bus configuration problems have been addressed in other contexts and a variety of solutions proposed.\nFor instance, U.S. Pat. No. 5,550,990 discusses physical partitioning of logically continuous buses. This document is directed to the SCSI (Small Computer System Interface) bus architecture, and suggests partitioning the bus into two or more physical entities which to the computer appear as one logical entity. This allows addressing problems potentially arising because of the scope of the architecture to be resolved; one example of such problems is excessive signal degradation due to use of signal rates which although allowed by the architecture are inappropriate for a particular bus loading. The solution disclosed in this document is to provide on the bus an adapter; instead of ensuring physical continuity of the bus, the adapter separates the bus into two bus partitions. This makes it possible, e. g. to operate the two partitions of the bus at different speeds, or to increase the number of devices connected to the bus. Where the speed has to be determined, a negotiation between the adapter and the devices connected to the bus is carried out at the time the adapter is initialised.\nU.S. Pat. No. 5,870,571 discusses automatic control of data transfer rates over a computer bus; this document is particularly directed to UltraSCSI buses. This document suggests detecting whether a SCSI external device is connected to the bus, and if this is the case, inhibiting the host adapter in order to reduce the data transfer rates to SCSI rate; otherwise, if no external SCSI device is detected, the UltraSCSI rate may be used, and the host adapter is not inhibited. In this document, the adapter polls the devices connected to the bus at initialisation, in order to know the transfer rate at which they may operate. Note that the operation of devices connected to the bus is not modified, since the host adapter only is inhibited.\nU.S. Pat. No. 5,237,690 discusses configuration at boot of IBM PS/2 personal computers. These computers provide a POS (programmable option select) for defining or providing settings for the assignment of system resources to a system board and various adapters. In order to avoid having to reconfigure the computer each time an adapter is added, removed or changed, this document suggests testing at boot of the computer whether any adapter was added, removed or changed; if this is the case, the adapters that were altered are disabled, and the computer is operated with all other adapters.\nU.S. Pat. No. 5,797,032 discusses a bus for connecting extension cards to a data processing system, and more particularly and ISA or EISA bus. For addressing the problem of collisions between the addresses of the different cards, this document suggests enabling all cards one at a time, for testing the addresses to which they respond. The cards that generate collisions are then disabled, and a message is displayed on a monitor for indicating to the user which cards were disabled.\nThe configuration constraints with which these two latter documents are concerned are logical-layer constraints and to resolve associated configuration problems the systems described rely on the buses concerned operating correctly at the physical level.\nAccording to the invention, there is provided a bus system comprising a controller; a high speed data transfer bus, the data transfer bus being subject to one or more inherent physical-layer configuration constraints for proper electrical operation; and a separate control bus, said control bus and said data transfer bus connecting the controller and the, or each, device connected thereto, wherein the controller is arranged to communicate with devices using the control bus in order to verify whether or not one or more of the physical-layer configuration constraints are satisfied and, if such configuration constraints are not satisfied, to modify using control signals transmitted on the control bus the operation of at least some of the devices in order to bring the data transfer bus to an operable condition.\nPreferably, if the configuration constraints are not satisfied, the controller is arranged to disable at least some of the devices using control signals transmitted on the control bus in order to bring the data transfer bus to an operable condition. The disabled devices may be the devices furthest from the controller on the data transfer bus. The controller may also be arranged to disable all devices connected to the bus, except one to five devices. The controller may also be arranged to set a stored indicator indicative of a error condition.\nIn one embodiment of the invention, the physical-layer constraints comprise a constraint on the number of devices connected to the bus.\nThe invention also provides a computer comprising such a bus.\nThe invention further relates to a process for bringing a data transfer bus to an operable condition in a bus system comprising a controller; a high speed data transfer bus, the data transfer bus being subject to one or more inherent physical-layer configuration constraints for proper electrical operation; and a separate control bus, said control bus and said data transfer bus connecting the controller and the, or each, device connected thereto. The process comprises the steps of\ncommunicating with devices using the control bus in order to verify whether or not one or more of the physical-layer configuration constraints are satisfied and,\nif such configuration constraints are not satisfied, to modifying the operation of at least some of the devices using control signals transmitted on the control bus.\nThe step of modifying may comprise disabling at least some of the devices using control signals transmitted on the control bus, and for instance, disabling devices furthest from the controller on the data transfer bus. The step of modifying may also comprises disabling all devices connect the bus, except one to five devices. The process may also comprise, if said configuration constraints are not satisfied, setting a stored indicator indicative of a error condition.\nIn one embodiment of the process, the physical-layer constraints comprise a constraint on the number of devices connected to the bus.\nThe invention also provides a computer program product for a computer with a bus system comprising a controller; a high speed data transfer bus, the data transfer bus being subject to one or more inherent physical-layer configuration constraints for proper electrical operation; and a separate control bus, said control bus and said data transfer bus connecting the controller and the, or each, device connected thereto. The computer program product comprises a computer readable medium having thereon:\ncomputer program code means, when said program is loaded, to make the controller communicate with devices using the control bus in order to verify whether or not one or more of the physical-layer configuration constraints are satisfied and,\nif such configuration constraints are not satisfied, to make the controller modify the operation of at least some of the devices using control signals transmitted on the control bus.\nPreferably, if such configuration constraints are not satisfied, the computer program code means make the controller disable at least some of the devices using control signals transmitted on the control bus. The disabled devices may be the devices furthest from the controller on the data transfer bus. The computer program code means may also make the controller disable all devices connected to the bus, except one to five devices. In another embodiment, the computer program code means set a stored indicator indicative of a error condition.\nThe physical layer constraints may comprise a constraint on the number of devices connected to the bus.\nThe invention provides a solution to the above described problem. The invention allows a computer at least to boot, even if the bus is improperly configured; this makes it possible to display a message to the user, so that he may address the problem. For Rambus, the mechanical configuration of the bus makes it possible to violate the bus specification by connecting an excessive number of devices on the bus.\nIn the case of the Rambus system, the limitation in the number of modules and devices connected to the bus is thought to be due to the sensitivity of the high speed signalling used on the Rambus Channel (the RSL or Rambus Signalling Levels) to the number of loads.\nIn consequence, a number of loads higher than the highest allowable number has no impact on the control bus, and does not affect the RSL signals for a few devices close to the controller. This makes it possible to disable some devices and to allow the bus to operate in a degraded operation mode. This mode is sufficient for booting a computer, and for allowing a warning to be communicated to the user, e. g. by displaying a message, so that they may reduce the number of devices on the bus and fix the problem.\nThe invention is however not limited to such a problem in the number of devices, but also can be applied in order to resolve other types of improper configuration; for instance, the invention could be applied if the bus can comprise different types of attachable devices, e.g. devices operating at different speeds, or devices requiring a special controller. It could also be applied for solving problems such as the mechanical length of a bus."}
{"text": "The present invention relates to a method of selecting transmission channel in a time division multiple access protocol and to a communication system using the method of the invention.\nIn the state of the art, for time division multiple access protocols, methods have already been proposed for selecting a transmission channel from a plurality of transmission channels as a function of the transmission conditions measured on each of those channels. These methods apply in particular to the case of a radio transmission between a base station and a set of mobile telephones, the quality of transmission of each channel linking the base station to a mobile telephone varying independently according to the mobility of the telephone and the change in its radio environment.\nThus a method described in document U.S. Pat. No. B1-6,449,490 consists in:                receiving for each channel a value representing the transmission rate currently achievable on that channel;        determining for each channel the transmission rate achieved on that channel during a time window;        selecting the channel for which the ratio of transmission rate currently achievable to the transmission rate achieved during a time window is the highest.        \nThe value of this method relative to a selection method that does not take account of the variations in the quality of transmission of the channels is to increase the transmission rate of each channel by selecting the one on which the transmission conditions are the most favourable. The selection is made on the basis of the ratio of the transmission rate achievable to the transmission rate achieved to provide relatively equitable access to the transmission resource. The drawback of this method is that a channel on which the achieved rate would be limited for a reason other than the sharing of the transmission resource (a constraint linked to the reception capability of the mobile telephone for example) could be selected systematically, even when its transmission conditions are bad.\nAnother method aiming to take advantage of the independent variations in the transmission quality of the channels and described in the same document consists in:                receiving for each channel a periodic indication of the transmission quality of that channel;        computing for each channel the average transmission quality of that channel;        selecting the channel for which the ratio of the last transmission quality indication to the average transmission quality is the highest.        \nWhen applying this teaching of the art, it emerges that certain channels would be allocated a much larger portion of the transmission resource than others. In fact, the distribution of the transmission resource depends on the distribution of the random variations of the transmission conditions of the set of channels. In addition, this distribution of the resource and the transmission quality of each resultant selected channel are extremely difficult to evaluate when the distributions of random variations of the transmission conditions of the channels differ, which is the case in practice."}
{"text": "1. Field of the Invention\nThe present invention relates to an anti-skid controlling apparatus which controls the braking force applied to the wheels of a vehicle so that the wheels are not locked when braked.\n2. Description of the Prior Art\nWhen a running vehicle is abruptly braked and its wheels are locked thereby, the vehicle's purchase on the road surface may become unstable depending on the road surface condition. To prevent the wheels from getting locked in such a situation, the anti-skid controlling apparatus has been conventionally used on board the vehicle. The controlling apparatus controls the brake fluid pressure in the wheel cylinders attached to the wheels so as to optimally decrease or increase the braking force applied to the wheels.\nAs the brake fluid pressure in the wheel cylinders is raised, the wheel velocity abruptly drops just before the friction coefficient .mu. of the wheels reaches its maximum. This characteristic is utilized by prior art anti-skid controlling apparatus that control the brake fluid pressure in keeping with changes in the wheel velocity and wheel acceleration. The braking force is controlled in this manner so that the slip ratio of the wheels settles at around 20% where a maximum friction coefficient is acquired.\nThe representative prior art anti-skid controlling apparatus works as follows. The rotating speed of each of the wheels, i.e., the wheel velocity, is first detected. The vehicle velocity is deducted from the wheel velocity. In turn, the vehicle velocity is used to obtain a reference velocity for comparison with the velocity of the individual wheels. The result of the comparison is used to control the brake fluid pressure in the wheel cylinder of each wheel. Japanese Patent Laid-open No. 64-52568 discloses one such anti-skid controlling apparatus. This apparatus distinguishes a \"pressure decreasing region\" from a \"pressure increasing region\" in its pressure control setup depending on the slip ratio and wheel acceleration of the vehicle on which it is mounted, each region being variable as required. The disclosure also includes an apparatus capable of lowering the memory size required for skid control.\nOne disadvantages of the prior art above is that simply separating the pressure decreasing region from the pressure increasing region fails to address the diverse control conditions that vary depending on many factors: different pressure changing characteristics of the pressure controlling means for controlling the brake fluid pressure on different vehicles; different tires; different characteristics of the braking system as a whole from vehicle to vehicle; and different friction coefficients relative to different road surfaces. The above Japanese Patent Laid-open No. 64-52568 discloses that in order to control brake fluid pressure decreases, the pressure decreasing action and pressure holding action are alternated by pulse-based pressure decreasing controls. In this setup, the ratio of pressure decreasing time to pressure holding time is predetermined for the pressure decreasing region. This results in another disadvantage: predestined difficulty in flexibly achieving optimal pressure control."}
{"text": "Light-emitting elements, such as light-emitting diodes, have recently come to be used widely as, for example, a source of light for electroluminescent display panels, or a source of backlight for cellular phones or personal computers, owing to the improved brightness of light thereby emitted. These light-emitting elements are used in packages in which they are stored.\nConventionally, as for a light-emitting element storing package, a light-emitting element storing package disclosed in Patent Document 1 (Japanese Patent Provisional Publication No. 2002-232017) is proposed.\nThe light-emitting element storing package 100 disclosed in the Patent Document 1, as shown in FIG. 6, has an insulating substrate 111, a reflector frame 112 disposed along the outer periphery of the substrate 111. Furthermore, in the package 100, a light-emitting element storing concave portion 119 is defined between the substrate 111 and the reflector frame 112.\nAt the position of the concave portion 119 of the insulating substrate 111, i.e. at the just under the position of the concave portion 119, a wiring through-hole (so called “via-hole”) 116, which is penetrated vertically through the substrate 111, is provided.\nThe wiring through-hole 116 is filled with a conductive member 117 and wiring pattern layers 114 and 115, which are respectively formed on the upper and lower surfaces of the substrate 111, are electrically connected to each other through the conductive member 117. As a result, the light-emitting element 118 mounted on the upper wiring pattern layer 114 may be electrically connected with an external electrical circuit.\nAfter the light-emitting element 118 is mounted on the bottom of the light-emitting element storing concave portion 119, the concave portion 119 is filled with a transparent resin, such as an epoxy resin so that the light-emitting element is sealed gas-tightly. As a result, a light-emitting device as a final product is produced.\nIn the light-emitting element storing package 100, a light-reflecting layer 113 is formed on the inner wall surface of the concave portion 119, i.e. on the inner wall surface of the reflector frame 112. The light emitted by the light-emitting element 118 is reflected by the light-reflecting layer 113 so that brightness of the light-emitting element 118 is substantially increased.\nAs a result, the light emitted by the light-emitting element 118 shines brightly by the action of this light-reflecting layer 113, so that an electroluminescent display panel may display bright and very clear characters and images.\nAs for the light-reflecting layer 113, for example, gold, silver or nickel, etc. is utilized. In addition, for example, a metallized metal layer is previously formed on the inner concave wall surface of the insulating substrate 111 and gold, silver or nickel, etc. is formed on the metallized metal layer by plating, so that the light-reflecting layer 113 is formed on the inner wall surface of the reflector frame 112.\nHowever, in such light-emitting element storing package 100 disclosed in the Patent Document 1, the light-reflecting layer 113 is formed on the whole inner wall surface of the reflector frame 112. Therefore, lower portion of the light-reflecting layer 113 and the wiring pattern layer 114 at the bottom of the light-emitting element storing concave portion 119 is short-circuited, so that there is a possibility of impairing the function of the light-emitting element 118.\nFor this reason, in a light-emitting element storing package 200 of Patent Document 2 (Japanese Patent Provisional Publication No. 2003-273405), as shown in FIG. 7, a light-reflecting layer 213 is formed on the inner wall surface of the light-emitting element storing concave portion 219 not to cover whole inner wall surface of the reflector frame 212. Moreover, in the package 200, the light-reflecting layer 213 is formed as it is spaced apart from the bottom of the light-emitting element storing concave portion 219.\nAs a result, the electrical short-circuit between the light-reflecting layer 213 and a wiring pattern layer 214 at the bottom of the light-emitting element storing concave portion 219 is prevented effectively.\nFurthermore, in the package 200, an insulating layer 220 is disposed between the light-emitting layer 213 and the wiring pattern layer 214 to ensure the prevention of their electrical short-circuiting.\nIn addition, in the light-emitting element storing package 200 disclosed in Patent Document 2, as shown in FIG. 7, there is no wiring through-hole in an insulating substrate 211, and a wiring pattern layer 201, which is provided between the reflector frame 212 and the substrate 211, is extended along the outer surface of the substrate 211. As a result, the wiring pattern layers 214 and 215, which are respectively formed on the upper and lower surfaces of the substrate 211, are electrically connected to each other.\nHowever, in the light-emitting element storing package 200 disclosed in this Patent Document 2, the light-reflecting layer 213 is formed as it is spaced apart from the bottom of the light-emitting element storing concave portion 219 and the insulating layer 220 is disposed therebetween. Therefore, the light emitted from the light-emitting element 218 is not reflected by the surface of the insulating layer 220 so that brightness is substantially lowered.\nFurthermore, according to both of Patent Documents 1 and 2, the light-reflecting layer is of, for example, gold, silver or nickel. However, gold and nickel have a light reflectance of as low as 50% or less for a wavelength of 400 nm or less. In addition, silver generally has a reflectance of as low as 50% or less for a wavelength of 350 nm or less, since it has absorption at a wavelength of 300 to 350 nm.\nHowever, a light-emitting element for a white LED usually emits blue or near-ultraviolet light and it is converted to various wavelengths by various fluorescent materials. Therefore, it is important for improving its brightness that a light emitted from a light-emitting element having a wavelength of 350 to 430 nm is efficiently reflected.\nTherefore, in Patent Document 3 (Japanese Patent Provisional Publication No. 8 (1996)-274378), as shown in FIG. 8, a light-reflecting layer is not formed and an insulating substrate 312 itself formed from a white sintered compact of aluminum oxide having relatively high light-reflecting property. As a result, a light emitted from a light-emitting element 318 is reflected by a light-reflecting surface 320.\nIn addition, in the light-emitting element storing package 300 disclosed in Patent Document 3, as shown in FIG. 8, a wiring through-hole 316 is formed. The wiring through-hole 316 is extending from the lower portion of a light-emitting element storing concave portion 319 of the insulating substrate 312 to the bottom surface of the insulating substrate 312 through an intermediate wiring portion 313 formed in the middle position of the insulating substrate 312.\nThe wiring through-hole 316 is filled with a conductive member 317 and wiring pattern layers 314 and 315, which are respectively formed on the upper and lower surfaces of the insulating substrate 312, are electrically connected to each other. Consequently, the light-emitting element 318, which is mounted on the upper wiring pattern layer 314, may be electrically connected with an external electrical circuit.\nHowever, the white sintered compact of aluminum oxide constituting the substrate 312 according to Patent Document 3 has a relatively low thermal conductivity in the order of about 20 W/m·K despite its reflectance of 50% or more for light having a wavelength of 350 to 400 nm. Therefore, the heat from the light-emitting element 318 is thermally accumulated in the light-emitting element storing package 300 and the accumulated heat can not radiate therefrom so that there is possibility of damaging the light-emitting element 318.\nMoreover, Patent Document 4 (Japanese Patent Provisional Publication No. 2004-152952) discloses a light-emitting element storing package in which a reflector frame is composed of white ceramics.\nHowever, in this Patent Document, the ceramics actually used for the reflector frame are SiO2—Al2O3—MgO—ZrO2—CaO ceramics having a specific composition, which is mainly consisting of alumina. As apparent from this Patent Document, the ceramic has 96.25% by weight of alumina content so that it has a low thermal conductivity.\nPatent Document 4 also states that, as a reflectance of 80% or higher for light having a wavelength of 400 to 700 nm with respect to aluminum nitride (AlN), the amount of Er2O3 in a sintered compact of aluminum nitride is preferably from 1 to 10% by weight based on the weight of the sintered compact.\nHowever, in Patent Document 4, there is no specific embodiment corresponding to this statement. In addition, in fact, it has been confirmed that no white sintered compact of aluminum nitride can be obtained by sintering a material having an Er2O3 content of 1 to 10% by weight (see Comparative Example 4 below).\nPatent Document 1 (Japanese Patent Provisional Publication No. 2002-232017)\nPatent Document 2 (Japanese Patent Provisional Publication No. 2003-273405)\nPatent Document 3 (Japanese Patent Provisional Publication No. 8 (1996)-274378)\nPatent Document 4 (Japanese Patent Provisional Publication No. 2004-152952)"}
{"text": "This specification relates to using synthetic descriptive text to rank search results.\nInternet search engines aim to identify resources, e.g., web pages, images, text documents, or multimedia content, that are relevant to a user's needs and to present information about the resources in a manner that is most useful to the user. Internet search engines generally return a set of search results, each of which identifies a resource, in response to a user submitted query."}
{"text": "1. Field of the Invention\nThe present invention relates to a differential amplifier circuit and, more particularly to, a differential amplifier circuit wherein the offset voltage is compensated for.\n2. Description of the Related Art\nFIG. 10 shows a conventional differential amplifier circuit comprising NPN transistors Q1 and Q2, current source A, and PNP transistors Q3 and Q4. Input signals V1 and V2 are supplied to the base of transistor Q1 and the base of transistor Q2, respectively. Current source A is connected between the common emitter of transistors Q1 and Q2 and the ground GND. PNP transistors Q3 and Q4 function as active loads of transistors Q1 and Q2, respectively. The offset voltage of this differential amplifier circuit will be discussed.\nThe bases of PNP transistors Q3 and Q4 are mutually connected, and the emitters of these transistors are coupled to power-supply Vcc terminal. Therefore, the base-emitter voltage Vbe3 of transistor Q3 is equal to the base-emitter voltage of transistor Q4. Hence: EQU (Ic3/Is3)=(Ic4/Is4)\nwhere Ic3 is the collector current of transistor Q3, Ic4 is the collector current of transistor Q4, Is3 is the reverse saturation current of transistor Q3, and Is4 is the reverse saturation current of transistor Q4.\nThe collector current Ic2 of transistor Q2 is given: EQU Ic2=Ic4=Ic3(Is4/Is3)\nThe collector current Ic1 of transistor Q1 is given: EQU Ic1=Ic3(1+2/.beta.p)\nwhere .beta.p is the current-amplification factor of PNP transistors Q3 and Q4.\nThus, the offset voltage Vos of the differential amplifier circuit shown in FIG. 10 is: ##EQU1## Here, let us assume that: ##EQU2## Then, the offset voltage Vos can be represented as follows: ##EQU3## Since .DELTA.Isp&lt;&lt;Isp, and .DELTA.Isn&lt;&lt;Isn, this equation reduces to: ##EQU4## Here is the relationship: EQU ln(1+x)=x-(x.sup.2 /2)+(x.sup.3 /3) . . .\nAnd .DELTA.Isp/Isp&lt;&lt;1, .DELTA.Isn/Isn&lt;&lt;1, and 2/.beta.p&lt;&lt;1. Therefore, the second term et seq. of the above equation can be neglected. Hence: ##EQU5## The values of .DELTA.Isp/Isp and .DELTA.Isn/Isn, both included in equation (1), are determined by the characteristics of the transistors used in the differential amplifier circuit. The third term in the right side, i.e., 2/.beta.p, is a factor which is determined by the structure of the circuit and the characteristics of the transistors used therein.\nAssuming that V.sub.T =26 mV, .DELTA.Isp/Isp=5%, .DELTA.Isn/Isn=5%, and .beta.p=10, then the worst value for Vos will be: EQU Vos=26 (0.05+0.05+0.1)=7.8 mV\nHence that portion of the offset voltage which is determined by the their term in the right side of equation (1) is 2.6 mV (=26.times.0.1).\nAs has been pointed, the third term in the right side of equation (1) is determined by the structure of the circuit and the characteristics of the transistors used in this circuit. When PNP transistor Q5 is added to the circuit shown in FIG. 10, with its base coupled to the collector of transistor Q3 its emitter connected to the common base of transistors Q3 and Q4, and its collector connected to the ground, as is illustrated in FIG. 11, then, equation (1) changes to: ##EQU6##\nWhen factors which are similar to those mentioned above, is applied, then: EQU Vos=26 (0.05+0.05+0.0182)=3.07 mV\nIn the circuit shown in FIG. 11, that portion of the offset voltage which is determined by the third term in the right side of equation (2) is 0.47 mV (=26.times.0.0182), far lower than the corresponding portion of the offset voltage of the differential amplifier circuit shown in FIG. 10. Nonetheless, the offset voltage (i.e., 2.6 mV) of the circuit is still too high."}
{"text": "The present invention relates to x-ray inspection of containers, and, more particularly, to x-ray inspection employing the detection of backscatter radiation by means of a plurality of backscatter detectors in order to derive information including spatial and material information with respect to contents of the containers.\nIt is desirable to be able to determine the presence of objects, such as contraband, weapons, or explosives, that have been concealed in an enclosure, such as luggage or a shipping container, or that are concealed behind a surface such as a wall. Additionally, it is desirable to obtain information regarding the geometrical and material characteristics of such objects. Conventional x-ray techniques provide measures either of attenuation, in the case of transmission techniques, or of scatter, in the case of scatter techniques.\nVarious methods of identifying a backscatter signal with a position within an illuminated enclosure employ scanned beams of x-rays, as described, for example, in U.S. Pat. Nos. 4,809,312 and 4,825,454 which are hereby incorporated herein by reference. In practice, the backscatter intensity may give only a crude measure of the atomic number of the object since the backscatter intensity is a function of several variables: the effective atomic number of the object; the object\"\"s geometry, including its distance from the x-ray source and the detectors; and the presence of material interposed between the object and the x-ray source/detector arrangement. It is desirable to obtain and process further scatter data so as to resolve a portion of the ambiguity inherent in backscatter measurements.\nIn accordance with one aspect of the invention, in a preferred embodiment, there is provided an inspection system for characterizing an object concealed by a concealing surface. The system has a source of penetrating radiation for emitting a beam that has an orientation and is incident upon the concealing surface at a plane of incidence. The source of penetrating radiation is characterized by a source position. The system also has a first scatter detector having a specified position with respect to the source position and beam orientation, for generating a first signal corresponding to penetrating radiation that has been scattered by the object. Additionally, the system has a second scatter detector having a field of view. The second scatter detector also has a specified position with respect to the source position and beam orientation and generates a second signal corresponding to penetrating radiation that has been scattered by the object. Finally, the system has a controller for determining an effective depth of the object with respect to the plane of incidence on the basis of at least the first and second signals."}
{"text": "The present invention relates generally to the field of can lids and specifically to lids for use on aluminum or metal beverage cans.\nThe applicant knows of no prior art which teaches the unique structure of his invention. U.S. Pat. No. 2,753,051 discloses a hinge on a pintle for sealing a container. Otherwise none of the structure of the present invention is shown. U.S. Pat. No. 3,372,832 (Yeater) discloses a removable cover for containers. Essentially it is a plastic cover having a pop down lid which locks into place by the force of the user pushing on the cap so that the projections 23 on the bottom side of the cap are pushed past the constricted middle section 16. The structure of the present invention is not disclosed. U.S. Pat. No. 4,331,255 (Fournier) discloses a two-part lid hinge. The two sections of this lid are connected by a hinge 16 which just may be a piece of plastic. The two sections are oriented so that the second section can be hooked over and on to the first so that the peripheral channel of the second section engages over the peripheral channel of the first section to provide a shallow space between the two sections bounded by a peripheral seal. Small openings are provided in a first section leaving a relatively large continuous imperforate area in that section. A removal tab is provided for forming a drinking opening in the second section, the tab being positioned so that it is disposed opposite the imperforate area in the first section. Again the structure of the present invention is not specifically disclosed. U.S. Pat. No. 3,977,559 (Lombardi) discloses a lid for a food container. It is a plastic cover that snaps over the top of a food container having a section of that lid that has been cut out so that it may be folded back along a hinge 22. The hinge does not have or disclose the present invention's structure nor does the cap itself disclose the structure of the present invention. U.S. Pat. No. 3,994,411 (Elfelt) discloses a lid for drink cups that includes a drinking flap of limited circumferential extent that may be selectively pivotally opened and closed. Such a drinking flap may be initially defined by frangible, i.e. breakable, lateral edges in the lid and may be held in its open position by the pull tab on the flap being inserted in a slit for a straw orifice. Essentially all this patent discloses, that is pertinent to the structure of the present invention, is a cap with hinges. The Elfelt structure does not appear to be re-usable. U.S. Pat. No. 4,284,200 (Bush) discloses a child resistant dispensing closure. Again, the structure of this invention is much different from the structure of the present invention and all that is really pertinent with this reference is the fact that there is a cap with a hinge. U.S. Pat. 4,361,250 (Foster) discloses a plastic container closure. This cap may be re-usable and has a flap or portion of the cover that flips open, similar in function to the present invention, in that the flap of Foster seals well. However, the structure of Foster is different from the structure of present invention. The hinged flap of Foster has hinge strips with depending pins formed along the sides of the flap that are integrally connected with the flap by tearable webs. After the flap has first been closed, the pins are anchored to the top of the closure and prevent the flap from opening during shipment of the container. Initial opening of the flap is effected by swinging the flap upwardly with a substantial force to tear the webs and separate the webs from the strips and the anchor pins. The torn webs provide visual indication that the flap has been opened. Accordingly, this is apparently a one use device or cap since once the webs have been torn the re-useability of the device is questionable. It is designed to be tamper resistant and tamper evident packaging. The structure of the present invention is not disclosed. U.S. Pat. 4,494,679 (Cleevely) discloses a thermoplastic container closure for dispensing solids. Again, this is another structure showing a hinged flap or flip cap. However, this flip cap does not appear to have a pintle hinge-type structure. U.S. Pat. No. 4,537,326 (Morehead) merely discloses a protector for a drink can opening. Specifically it is a device designed to attach to the portion of the lid where the pop tab or the pull tab is located so that it can be swung over or incorporated into the structure of the pop can lid in order to have a grating in place over the drinking opening and prevent the ability of insects like bees to enter into the container. Other than the slits the structure is completely different from the present invention. U.S. Pat. No. 4,582,214 (Dart) discloses a non-spill drink through lid. Slits to drink through are disclosed. No other structure of present invention is disclosed. U.S. Pat. No. 4,619,372 (McFarland) discloses a cap for a hot beverage cup. The cap is a disposable, removable closure cap for beverage containers and includes a depression permitting the beverage to be drunk while the cap remains in place on the container. The cap includes perforations in a depending wall located closely adjacent an inner wall of the container in order to limit the flow of beverage into the depression defined in the upper surface of the cap. A slit in the cap permits the aroma of the beverage to be enjoyed while the beverage is being drunk from the depression. The cap may be formed from sheet plastic material. This structure is completely different from the structure of the present invention, although it does disclose slits. It does not disclose the lip protection feature of the present structure although the well 28 of the McFarland device is designed to produce a somewhat similar effect. However the structure is completely different. U.S. Pat. No. 4,629,088 (Durgin) discloses a beverage container lid which includes a foldable flap which may be opened to allow the user to drink from a beverage container which is covered by the lid. A recess in the beverage lid is provided to receive the opened flap and to firmly secure the flap in its open position. The recess includes a pair of detentes on either side and an overhang at one end which cooperate to hold the flap firmly within the recess. The flap of course is also hinged. The structure of this cap is different than the structure of the cap which is the present invention, although it appears that recess 40 allows the flap to be flipped back and locked in place so that it is out of the way when a person drinks from the container. U.S. Pat. No. 4,949,865 (Turner) discloses a container lid with an integral stop. Essentially the structure of Turner is quite different from the present invention's structure and the main similarity is the fact that it discloses a hinged flap on the lid. The hinge is molded unitarily to the upper margin of the central support in the cover. Finally, U.S. Pat. No. 4,796,774 (Nabinger) also discloses a removable and re-sealable lid for a container but Nabinger's structure is much different from the present invention's structure."}
{"text": "This invention relates to novel hydratropic acid derivatives."}
{"text": "The polymerase chain reaction (PCR) is a method for increasing by many orders of magnitude the concentration of a specific nucleic acid sequence in a test sample. The PCR process is disclosed in U.S. Pat. Nos. 4,683,195; 4,683,202; and 4,965,188, each of which is incorporated herein by reference.\nIn PCR, a test sample believed to contain one or more targeted nucleic acid sequences is combined in a total volume of usually about 20 to 200 .mu.l with the following reagents: an aqueous buffer, pH 8-9 at room temperature, usually also containing approximately 0.05 M KCl; all four common nucleoside triphosphates (e.g., for DNA polymerase, the four common dNTPs: dATP, dTTP, dCTP, and dGTP) at concentrations of approximately 10.sup.-5 M to 10.sup.-3 M; a magnesium compound, usually MgCl.sub.2, usually at a concentration of about 1 to 5 mM; a polynucleotide polymerase, preferably a thermostable DNA polymerase, most preferably the DNA polymerase I from Thermus aquaticus (Taq polymerase and the Stoffel fragment of Taq polymerase are the subject of U.S. Pat. No. 4,889,818, incorporated herein by reference; the latter enzyme lacks the 5'.fwdarw.3' exonuclease activity of native Taq polymerase), usually at a concentration of 10.sup.-10 to 10.sup.-8 M; and single-stranded oligonucleotide primers, usually 15 to 30 nucleotides long and usually composed of deoxyribonucleotides, containing base sequences which are Watson-Crick complementary to sequences on both strands of the target nucleic acid sequence(s). Each primer usually is present at a concentration of 10.sup.-7 to 10.sup.-5 M; primers are synthesized by solid-phase methods well known in the art of nucleic acid chemistry.\nIn the simplest form, PCR requires two primers for each target sequence. These primers, when annealed to the opposing target strands, have their 3' ends directed toward one another's hybridization sites and separated by about 100 to 1,000 nucleotides (occasionally up to about 10,000 nucleotides). The polymerase catalyzes magnesium-dependent, template-directed extension of each primer from the 3' end of the primer, incorporating nucleoside monophosphates into the growing nucleic acid and releasing pyrophosphate.\nThis extension reaction continues until the polymerase reaches the 5' end of the template strand to which the extended primer was annealed, at which point the polymerase is free to bind to another primer-template duplex and catalyze extension of that primer molecule; the extension reaction also stops if the reaction mixture is heated to temperatures sufficient to separate the template from the extended primer before the enzyme has reached the 5' end of the template. After the enzyme has worked long enough to transform a large fraction of the primer-template duplexes into double-stranded nucleic acid, the latter can be denatured at high temperature, usually 90.degree. to 100.degree. C., to create two single-stranded polynucleotides, which, after cooling to a temperature where they can be annealed to new primer molecules, serve as templates for another round of enzyme-catalyzed primer extension. Because both DNA strands serve as template, each round of nucleic acid replication approximately doubles the concentration of the specific nucleic acid sequence defined at its ends by the two primer sequences. Therefore, the total concentration increase in the target nucleic acid sequence in a PCR amplification is by a factor of approximately 2.sup.n, where n is the number of completed thermal cycles between a high temperature where double-stranded DNA is denatured and a lower temperature or set of temperatures (40.degree. to 75.degree. C.) where primer-template annealing and primer extension occur.\nAlthough one can move PCR reaction tubes manually back and forth between thermostated baths in the two temperature ranges, PCR most commonly is performed in an automated temperature-controlled machine, known as a \"thermal cycler,\" in which a microprocessor is programmed to change the temperature of a heat-exchange block or bath containing reaction tubes back and forth among several specified temperatures for a specified number of cycles, holding at each temperature for a specified time, usually on the order of one-half to two minutes. Such a thermal cycler is commercially available from Perkin Elmer Cetus Instruments and described in the European Patent Publication No. 236,069 and U.S. patent application Ser. No. 670,545, filed Mar. 14, 1991, which is a continuation-in-part of Ser. No. 620,606, filed Nov. 29, 1990, both of which are incorporated herein by reference. The total cycle time is usually less than 10 minutes, and the total number of cycles is usually less than 40, so that a single, multi-cycle amplification, amplifying the targeted nucleic acid sequence 10.sup.5 to 10.sup.10 times, normally takes less than seven hours and often less than four hours.\nThe practical benefits of PCR nucleic acid amplification have been rapidly appreciated in the fields of genetics, molecular biology, cellular biology, clinical chemistry, forensic science, and analytical biochemistry, as described in the following review volumes and articles: Erlich (ed.), 1989, PCR Technology, Stockton Press (New York); Erlich et al. (eds.), 1989, Polymerase Chain Reaction, Cold Spring Harbor Press (Cold Spring Harbor, N.Y.); Innis et al., 1990, PCR Protocols, Academic Press (New York); and White et al, 1989, Trends in Genetics 5/6:185-189. PCR can replace a large fraction of molecular cloning and mutagenesis operations commonly performed in bacteria, having advantages of speed, simplicity, lower cost, and sometime increased safety. Furthermore, PCR permits the rapid and highly sensitive qualitative and even quantitative analysis of nucleic acid sequences, often with greatly increased safety because so much PCR product is made that nonisotopic detection modes suffice.\nDespite rapid and broad adoption of PCR by a range of biological and chemical disciplines, PCR has sometimes suffered from the occurrence of side reactions which interfere with amplification of the specific target sequence or sequences. Many amplifications yield non specific side products differing in size and sequence from the sequence targeted by the primers used. Sometimes nonspecificity is caused by mis-priming, where primers have been annealed to non-target sequences, also present in the nucleic acid of the test sample similar to the target sequence. Although the genomic DNA commonly contained in PCR test samples has customarily been thought to be completely double-stranded, procedures used to prepare DNA for amplification appear to render that DNA, to a significant extent, single-stranded. Single-stranded DNA is especially susceptible to mis-priming if mixed with a complete set of PCR reagents at ambient or sub-ambient temperatures. Many PCR reagents also yield primer dimers or oligomers, double-stranded side products containing the sequences of several molecules joined end-to-end, the yield of which correlates negatively with the yield of amplified target sequence.\nRecently several methodological modifications have improved PCR specificity and sensitivity significantly. In Hot Start.TM. PCR, complete mixing of PCR reagents and test sample is delayed until reactants have been heated to a temperature in the 50.degree. C.-80.degree. C. range, sufficient to minimize mis-priming and primer dimerization; thermal cycling is started immediately after mixing at elevated temperature. In manual Hot Start.TM. PCR, the operator heats the reaction tube, containing test sample and a subset of PCR reagents, to the elevated incubation temperature, opens each tube separately to add a small volume of liquid containing the missing reagent(s), and closes each tube before moving on to the next one. See Frohman et al., 1988, Proc. Natl. Acad. Sci. USA 85:8998-9002; Ward et al., 1989, Nature 341:544-546; Newton et al., 1989, Nucl. Acids Res. 17:2503-2516; and Faloona et al., Abstract 1019, 6th International Conference on AIDS, June 20-24, 1990, San Francisco, Calif. More recently, Hot Start.TM. PCR was rendered more convenient and precise by (1) replacement of the conventional mineral oil vapor barrier by a layer of wax melting in the 50.degree. C. to 80.degree. C. range, (2) assembly of reaction tubes such that before thermal cycling, PCR reactants are grouped into subsets separated by a solid wax layer, and (3) convective mixture of all reactants during the first heating step of thermal cycling after the solid wax melts into a lighter-than-water oil. Such wax-mediated, Hot Start.TM. PCR is the subject of U.S. patent application Ser. No. 481,501, filed Feb. 19, 1990, now abandoned in favor of continuation application U.S. Ser. No. 07/890,300, filed May 27, 1992, incorporated herein by reference.\nAlternatively, nonspecific amplified nucleic acid resulting from primer dimerization and mis-priming while completely mixed PCR reactants stand at room temperature before thermal cycling can be destroyed by an enzymatic restriction process described in PCT U.S. patent application Ser. No. 91/05210 filed Jul. 23, 1991, which published as PCT Patent Publication No. WO 92/01814 on Feb. 6, 1992, which is a continuation-in-part of U.S. Ser. No. 609,157, filed Nov. 2, 1990, now abandoned which is a continuation-in-part of U.S. Ser. No. 557,517, filed Jul. 24, 1990, now abandoned each of which is incorporated herein by reference. To perform such restriction, one of the conventional four dNTPs is replaced by a structural analogue which is incorporated into all amplified nucleic acid by the PCR polymerase. Also included in the reaction mixture is an enzyme which digests nucleic acid at (and only at) positions which contain the structural analogue; this enzyme must be active only at temperatures below about 50.degree. C., so that it does not damage amplified nucleic acid during thermal cycling at higher temperatures. Preferably the restriction enzyme is permanently inactivated during thermal cycling, so that it cannot damage amplified nucleic acid if the latter is stored for any significant period of time at room temperature after amplification and before analysis. The most practical restriction enzymes are glycosidases which cleave from the polynucleotide phosphodiester backbone the unconventional nucleic acid base introduced by the dNTP analogue. The resulting abasic sites experience cleavage of the polynucleotide phosphodiester backbone upon heating. This restriction process has been integrated practically with PCR by replacing dTTP with dUTP and by incorporating in the reaction mixture the enzyme uracil-N-glycosidase.\nA chemical variant of the Hot Start.TM. process incorporates into the PCR reagent mixture a single-stranded DNA binding protein (SSB) at a concentration sufficient to bind a significant fraction of the single-stranded DNA present before thermal cycling is started. This ssDNA comprises minimally the primers, the concentrations of which are well known by the operator, and may also include slight or considerable amounts of the test sample DNA, depending on whether the latter has been prepared in a way which might denature it. During thermal cycling, the binding of the SSB to primers and single-stranded template strands formed by PCR product denaturation must be weak enough not to interfere with primer-template annealing and enzymatic primer extension. Before thermal cycling, while reactants stand together at room temperature, SSB binding to the primers and any single-stranded regions of test sample DNA must be strong enough to block mis-priming and primer dimerization. Two heavily studied SSBs (Chase and Williams, 1986, Ann. Rev. Biochem. 55:103-136) are commercially available and have been used with PCR: gene 32 protein from the bacteriophage T4 and the 19 kilodalton SSB from E. coli (19 kda is the subunit size; the normal active species is a tetramer). SSB is the major active ingredient of Perfect Match.TM. polymerase enhancer, a mixture of E. coli SSB and bovine serum albumin sold by Stratagene (San Diego, Calif.) for the purpose of increasing PCR specificity and yield. Bacteriophage gene 32 protein has been included in PCR mixtures to improve amplification of long targets (Schwarz et al., 1990, Nucl. Acids Res. 18:1079) and to relieve polymerase inhibition by blood in the test sample (Panaccio and Lew, 1991, Nucl. Acids Res. 19:1151). However, essentially all organisms possess SSBs with compositions unique to each organism. Other SSBs which have been characterized biochemically include one from a filamentous bacteriophage (Brayer and McPherson, 1984, Biochemistry. 23:340-349), a family of sequence-homologous proteins from plant virus (Saito et al., 1988, Virology 167:653-656, and Citovsky etal., 1990, Cell 60:637-647), and one from Agrobacterium tumefaciens (Citovsky et al., 1989, Proc. Natl. Acad. Sci. USA 86:1193-1197). SSBs possess enough structural similarity to suggest that DNA binding is associated with a consensus structure of alternating aromatic amino acids (phenylalanine, tyrosine, and tryptophan) and charged amino acids (glutamate, aspartate, lysine, and arginine) (Prasad and Chiu, 1987, J. Mol. Biol. 193:579-584) such that artificial polypeptides might be created which function as well as the biological SSBs in improving PCR specificity and yield. In addition, enough is known about SSB structure and function to suggest ways to improve function by genetic engineering.\nAlthough the three basic tactics of PCR specificity enhancement (Hot Start.TM. methods, amplified DNA restriction, and SSB addition to the reaction mixture) each can serve alone to improve specific amplification, combinations of the three approaches may have special benefits. For example, whereas Hot Start.TM. methods block only that nonspecificity resulting from reactant incubation at ambient temperature before cycling is started, SSB s may reduce mis-priming which arises during thermal cycling. On the other hand, SSB used without a manual or wax-mediated Hot Start.TM. process occasionally will trigger massive primer dimerization which interferes with specific amplification. The combination of the two methods optimally reduces mis-priming and primer dimerization.\nThe preceding background art has dealt with conventional PCR, wherein test sample nucleic acids are extracted from a biological source in a way which destroys target sequence association with individual cells or subcellular structures. So-called in situ nucleic acid hybridization methods have evolved to detect target sequences in the cells or organelles where they originated (for a review of the field, see Nagai et 1987, Intl. J. Gyn. Path. 6:366-379). Typically, in situ hybridization entails (1) preparation of a histochemical section or cytochemical smear, chemically fixed (e.g., with formaldehyde) to stabilize proteinaceous subcellular structures and attached to a microscope slide, (2) chemical denaturation of the nucleic acid in the cellular preparation, (3) annealing of a tagged nucleic acid probe to a complementary target sequence in the denatured cellular DNA, and (4) localized detection of the tag annealed to target, usually by microscopic examination of immobilized nonisotopic (absorbance or fluorescence staining) or isotopic (autoradiographic) signals directly or indirectly generated by the probe tag. However, conventional in situ hybridization is not very sensitive, generally requiring tens to hundreds of copies of the target nucleic acid per cell in order to score the presence of target sequence in that cell.\nRecently, the sensitivity enhancement associated with PCR amplification of target sequence has been combined with the target localization of in situ hybridization to create in situ PCR, wherein PCR is performed within chemically fixed cells, before (Haase et al., 1990, Proc. Natl. Acad. Sci. USA 87:4971-4975, incorporated herein by reference) or after (Nuovo et al., 1991, Amer. J. Pathol. in press, incorporated herein by reference) the fixed cells have been attached to a microscope slide; the amplified nucleic acid is located by microscopic examination of autoradiographs following isotopically tagged probing (Haase et al., supra) or stained patterns directly deposited on the microscope slide following enzyme-linked detection of biotin-tagged probes (Nuovo et al., supra). The cells may be suspended (Haase et al., Supra) or may be part of a tissue section (Nuovo et al., supra) during in situ amplification.\nIn situ PCR requires a delicate balance between two opposite requirements of PCR in a cellular preparation: the cell and subcellular (e.g., nuclear) membranes must be permeabilized sufficiently to allow externally applied PCR reagents to reach the target nucleic acid, yet must remain sufficiently intact and nonporous to retard diffusion of amplified nucleic acid out of the cells or subcellular compartments where it is made. In addition, the amplified nucleic acid must be sufficiently concentrated within its compartment to give a microscopically visible signal, yet remain sufficiently dilute that it does not reanneal between the denaturation and probe-annealing steps. Haase et al., supra, relied on paraformaldehyde fixation of cells to have created sufficient but not excessive permeability. Nuovo et al., supra, also employed a single, commercially available, proteinase treatment to improve permeability.\nBoth Haase et al., supra, and Nuovo et al., supra, used a series of PCR primer pairs to specify a series of overlapping target sequences within the genome of the targeted organism to improve retention of amplified target nucleic acid within the cellular compartment where it was made. The resulting PCR product was expected to be so large (greater than 1,000 base pairs) that its diffusion from site of origin should be greatly retarded. However, the use of multiple primer pairs severely reduces the practicality of in situ PCR, not just because of the expense associated with producing so many synthetic oligonucleotides, but even more seriously because many PCR target organisms, especially pathogenic virus, are so genetically plastic that it is hard to find even a few short sequences which are sufficiently invariant to make good primer and probe sites. Other important target sequences, such as activated oncogenes, inactivated tumor suppressor genes, and oncogenic chromosomal translocations, involve somatic point mutations and chromosomal rearrangements which can be distinguished from the parental sequence if relatively short PCR products are amplified from single primer pairs. Multiple primer pairs and long structures would frustrate attainment of the specificity often needed to distinguish cancerous cells from their normal neighbors. Multiple primer pairs jeopardize PCR in a different way as well; they promote primer dimerization and mis-priming, reducing sensitivity and specificity and increasing the likelihood of false-negative results because nonspecific amplification radically reduces the yield of amplified target sequence. Reinforcing the tendency of multiple primer pairs to enhance nonspecific amplification are the rather high primer concentrations preferred for in situ PCR (Nuovo et al., supra).\nOne useful variant of conventional PCR detects target RNA sequences in test samples by creating complementary DNA (cDNA) sequences with the catalytic mediation of added reverse transcriptase; the cDNA then is subjected to standard PCR amplification (Kawasaki et al., 1988, Proc. Natl. Acad. Sci. USA 85(15):5698, and Rappolee et al., 1989, J. Cell. Biochem. 39:1-11). Recently, such RNA PCR has been streamlined by using a thermostable DNA polymerase which, depending on exact chemical conditions, also shows strong reverse transcriptase activity. This enzyme and its optimized application to RNA PCR are subject of PCT U.S. patent application Ser. No. US90/07641, filed Dec. 21, 1990, incorporated herein by reference. Adaptation of in situ PCR to RNA targets will realize the full potential of the method to differentiate among neighboring cells in a histochemical or cytochemical preparation with respect to somatic mutation, pathogenic infection, oncogenic transformation, immune competence and specificity, state of differentiation, developmental origin, genetic mosaicism, cytokine expression, and other characteristics useful for understanding both normal and pathological conditions in eukaryotic organisms.\nThe present invention increases the convenience, sensitivity, and specificity of in situ PCR, also eliminating any need for multiple primer pairs to detect a single target sequence. In doing so, it also allows in situ PCR to discriminate among alleles and increases the practicality of in situ PCR analysis of RNA targets."}
{"text": "The invention relates to a shaft for manually guided lift vehicles and lift loaders equipped with a driving mechanism which has a steering head at its free end that contains gripping handles and operating switches for the driving and lifting operation of the vehicle.\nA steering head on a shaft of a lift vehicle is known which comprises operating gripping handles directly connected with the stem of the shaft as well as having actuating control switches for the lift vehicle. Gripping handles are supported at a stem of the steering head by means of bows. The steering heads, during operation of the lift vehicles are subjected to multiple stress as a result of improper operating and handling which may damage parts of the steering head making the lift vehicle inoperative.\nIt is an object of the invention to provide an improved steering head which has an easily exchangeable design.\nAccording to the invention, this object is achieved by having the steering head comprising two cup elements which are detachably connected with one another. One cup element is fixedly connected with the shaft and is constructed as a supporting part for the second cup element, which contains the gripping handles as well as the operating switches.\nIt is advantageous if the supporting part is provided with at least two supporting noses which are arranged at a distance from each other and at opposite sides of a stem of the shaft. Corresponding supporting noses are provided on the second cup element, disposed opposite the forementioned supporting noses in a contacting manner. Fastening screws for connecting the two cup elements with one another are located in the noses.\nIt is also advantageous if the second cup element has an arcuately shaped cross-section and a bent connecting leg with a front edge which is provided with a surrounding shoulder into which a front edge of the supporting part overlaps from the outside. By such a construction the two cups are form-lockingly held together such that the joint that locks the outer surfaces of the two cup elements extends in a flush manner at the connecting area. To provide better support is preferable if the two cup elements each have receiving bores for the fastening screws that have axes that extend at a distance from the switch and from the gripping handle.\nPlug in connections can be used to connect electrical connecting lines (from the switches on the second cup element) to the feeding lines in the supporting part.\nThe principal advantage achieved by the invention is, that as a result of its two-piece construction, the steering head is easily and rapidly exchangeable when the part of the steering head receiving the operating elements is damaged. This limits the time that the vehicle cannot be used and does not require the vehicle to be taken out of service.\nIn order to achieve this rapid exchanging of the one element of the steering head in a simple manner, the head is formed of two parts connected by way of two screws and with the supporting part at the shaft stem. The fastening screws can be reached in a straight line by means of a screwdriver without any impairment by any parts of the steering head.\nBy means of the special development of the two cup elements of the steering head with corresponding supporting noses and mutually engaging connecting edges, a sturdy structural unit is provided. Here the cup element (connected with the stem) essentially absorbs the forces occurring at the steering head.\nSo that a fast and problem-free disconnecting and connecting of the electrical lines takes place, the electrical feeding lines are connected with the other lines leading to the operating switches through one or more electrical plugs.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "(1) Field of the Invention\nThe present invention relates to a mobile phone, and especially to processing performed by a mobile phone upon writing a PRL (Preferred Roaming List).\n(2) Description of the Related Art\nSome telecommunications carriers have a mutual agreement to provide a service called roaming, which is the extending of connectivity in another carrier's service area. The roaming service allows users to use their mobile phones via a communication system (i.e. base station) of a carrier other than the carrier with which the users are registered. A list called PRL indicates which of available communication systems the mobile phone should use by priority, and also indicates other information regarding the respective communication systems.\nWith the use of PRL, the users are allowed to set up a selection condition, such as a communication mode, that needs to be satisfied by a communication system to be selectable. The mobile phone selects a communication system to be used, based on the PRL stored therein and the selection condition set by the user. That is, by setting up a selection condition specifying a specific communication mode, the user can ensure that the mobile phone conducts communications with a communication system in the desired communication mode.\nFor example, the user may set up such a selection condition that the mobile phone does not use any communication system available via a roaming service. Use of a roaming communication system generally incurs extra charges. By setting a selection condition for not using a roaming service, the user can save such extra charges.\nThe PRL stored in a mobile phone needs to be updated in the case where the telecommunications carriers build a new base station or remove an existing base station. In such a case, information contained in the PRL may no longer be consistent with the current state of communication systems. In order to maintain the accuracy of the PRL, each time a change is made to the communication systems, the mobile phone newly acquires a PRL reflecting the change, and selects a communication system to use based on the newly acquired PRL. A method for updating data stored in a mobile device from an external source is disclosed in JP Patent Application Publication No. 2001-217847.\nIt should be noted, however, that the newly acquired PRL may not contain any communication system satisfying the selection condition set by the user.\nSuppose, for example, the user specifies an analog communication mode as a selection condition, whereas all the communication systems contained in the newly acquired PRL are capable of digital communications and not capable of analog communications.\nIn the case where, as described above, the PRL does not contain any communication system satisfying the user-set selection condition, the mobile phone cannot select any communication system. As a consequence, the mobile phone will be unable to carry out communications.\nIn order to avoid such an undesirable occurrence, upon acquisition of a new PRL, a state-of-the-art mobile phone automatically resets the selection condition set by the user, so that all the communication systems are made selectable. This arrangement serves to prevent that the mobile phone will be unable to conduct communications.\nIt should be noted, however, according to the above arrangement, each time a new PRL is acquired, the user-set selection condition is automatically discarded and changed to a predetermined state. This change of the selection condition takes place without informing the user or without user's consent.\nAccording to OTAPA (Over-the-air Parameter Administration), the mobile phone automatically acquires a new PRL from an external source without requiring any user operation. Thus, the PRL acquisition and the change of selection condition take place while the user is not aware of it. As a result, it is possible that while the user believes that the mobile phone sets up communications using a communication system satisfying the selection condition set by the user, a communication system actually used is the one not desired by the user."}
{"text": "The inventor of the present invention has recognized the success of semiconductor devices has been mainly driven by an intensive transistor down-scaling process. However, as field effect transistors (FETs) approach sizes less than 100 nm, physical problems such as short channel effect begin to hinder proper device operation. For transistor based memories, such as those commonly known as Flash memories, other performance degradations or problems may occur as device sizes shrink. With Flash memories, a high voltage is usually required for programming of such memories, however, as device sizes shrink, the high programming voltage can result in dielectric breakdown and other problems. Similar problems can occur with other types of non-volatile memory devices other than Flash memories.\nThe inventor of the present invention recognizes that many other types of non-volatile random access memory (RAM) devices have been explored as next generation memory devices, such as: ferroelectric RAM (Fe RAM); magneto-resistive RAM (MRAM); organic RAM (ORAM); phase change RAM (PCRAM); and others.\nA common drawback with these memory devices include that they often require new materials that are incompatible with typical CMOS manufacturing. As an example of this, Organic RAM or ORAM requires organic chemicals that are currently incompatible with large volume silicon-based fabrication techniques and foundries. As another example of this, Fe-RAM and MRAM devices typically require materials using a high temperature anneal step, and thus such devices cannot be normally be incorporated with large volume silicon-based fabrication techniques.\nAdditional drawbacks with these devices include that such memory cells often lack one or more key attributes required of non-volatile memories. As an example of this, Fe-RAM and MRAM devices typically have fast switching (e.g. “0” to “1”) characteristics and good programming endurance, however, such memory cells are difficult to scale to small sizes. In another example of this, for ORAM devices reliability of such memories is often poor. As yet another example of this, switching of PCRAM devices typically includes Joules heating and undesirably require high power consumption.\nFrom the above, improved semiconductor memory devices that can scale to smaller dimensions with reduced drawbacks are therefore desirable."}
{"text": "Various devices and systems are known for use in harvesting solar energy by the use of photovoltaic cells. These include slat concentrators, which are photovoltaic devices generally comprising a series of parallel trough-shaped off axis parabolic reflectors to concentrate sunlight on photovoltaic receptors mounted on respective adjacent reflectors. The reflectors are typically automatically actuated to track the sun in order to ensure accurate reflection and concentration of solar radiation on the photovoltaic receptors.\nThe photovoltaic receptors forming part of such concentrators have a limited lifespan and the photovoltaic devices therefore require periodic removal and replacement of the photovoltaic receptors. There is a relationship between the operating temperatures of the photovoltaic receptors and their lifespan. Additionally, a photovoltaic receptor generally displays higher efficiency at lower temperatures."}
{"text": "Grid computing is a method of harnessing the power of many computational resources in a network. Grid computing is a distributed computer infrastructure involving large-scale sharing applications and/or high performance processing and network bandwidth.\nReferring to FIG. 1, a grid computing environment comprising a plurality of heterogeneous computing nodes distributed across multiple administrative domains. A virtual organization (e.g., VO-X, VO-Y, and VO-Z) may include nodes from different domains. For example, VO-Z includes nodes from each of Administrative Domains 110, 120, and 130.\nA node (hereinafter also referred to as a computing resource) may be a member of several virtual organizations. An end-user (e.g., USER M, USER N, or USER O) may need to access remote nodes either in the same administrative domain, or across domains. For example, User O may access VO-Y from within Domain 120, or User N within Domain 130 may access the same node from outside of Domain 120. Similarly, User M of Domain 130 may access both VO-Y and VO-Z through a single node in Domain 110.\nTraditionally, grid computing has provided for the execution of batch jobs in the scientific and academic community. Batch execution on a grid computing environment requires authentication, authorization, resource access, resource discovery, and other services. In support of batch processing of jobs on a grid computing environment, protocols, services, application programming interfaces, and software development kits have been developed. The conventional method and system are not particularly suited for interactive grid computing sessions.\nIn a grid computing environment, security can be a concern. As a result, a grid computing service, providing access to desktops for use, may reside behind a firewall. In addition, a user that is accessing a grid computing resource, may reside behind another firewall. Conventional methods for providing grid computing resources cannot traverse firewalls and provide a secure connection simultaneously.\nPort forwarding can be used to traverse firewalls between organizations. In this solution a third party host is needed between firewalls which is accessible from both organizations. Both organizations need to connect to this bridge host to open a tunnel.\nThis solution may be fine for academic settings, but in commercial environments, the client may not want to go over an intermediate bridge for connecting to the ASP. Also, bridge host needs to be especially set up in order to do port forwarding.\nA grid computing infrastructure that can hide the complexities of resource management from the user and that can provide secure graphical and multimedia interactive sessions is needed."}
{"text": "The present invention relates to a process for the early detection of damage to machine parts, particularly rolling element bearings, in which process forced vibrations are detected by means of a vibration pickup.\nBritish Patent No. 1,367,773 discloses a process for the early detection of damage to machine parts, particularly rolling element bearings, in which process signals from forced vibrations stemming from the machine are detected by means of a vibration pickup. Peak values of the forced vibrations always occur whenever a flaw in an otherwise undamaged surface rolls over a second undamaged surface. The picked up signals are transmitted to a peak value detector and to a mean value calculator. By means of a divider, the ratio of the output signal of the peak value detector to the output signal of the mean value calculator is formed and transmitted to a display unit. If the ratio exceeds a preselected limit, a signal is activated. Surface defects can be detected by means of this process, as long as they are still low in number and extent. However, in the case of this process, access to previously measured results is necessary to set the limit values of the ratio.\nA process of the type mentioned above, in which forced vibrations of the machine are detected at the machine by means of vibration pickups, is known from West German Patent No. 3,245,505. For this purpose, vibration pickups are used to measure a first structure-borne noise parameter in the high-frequency range and a second structure-borne noise parameter in the low-frequency range of the frequency spectrum, and both measured values are compared by means of a comparator. On the basis of the measured values of the first structure-born noise parameter and of the second body-noise parameter, an auto-performance spectrum G.sub.11 or G.sub.22, respectively, is formed for each of these parameters. A cross-performance spectrum G.sub.12 is also formed calculated on the basis of both structure-borne noise parameters. A coherence function is formed from these values and a rise in the coherence function serves as a signal for the early detection of damage. Although this process makes it possible to carry out early detection of damage independent of the previous determinations of reference values, this process is not suitable for estimating the remaining service life of the machine part."}
{"text": "Online auctions delivered to users via the World Wide Web have become a popular activity for users interested in purchasing goods and services at the least possible amount online. Many online auctions allow a person to bid on an item in an attempt to obtain a winning bid and an associated opportunity to purchase the auction item at the winning bid price. Known online auction websites allow several users to compete with other online auction participants by submitting bids on a particular item until a predefined period of time has elapsed, and the highest bidder is determined.\nAt least some known online auction websites require users to purchase a bid unit that represents an opportunity to submit a bid in an auction and to redeem a bid unit each time the user places a bid on an auction item. Once the user has redeemed all of the bid units, the user cannot participate in the auction until additional bid units are purchased or provided. By requiring a user to purchase bid units to participate in the auction the online auction receives revenue from each user participating in the auction.\nMany auction participants are attracted by the opportunity to win new auction items of high value at a low price, but may become frustrated with a requirement to purchase new bid units each time a user wishes to place a bid on an auction item. Auction participants may also become frustrated by purchasing bid units, participating in the auction, and not winning an auction. Likewise, online auction providers desire new auctioning opportunities to appeal to their auction participants, entice users to place additional bids, spend more time in the auction website, and have additional opportunities to benefit from winning new auctions. Accordingly, there is a continued need for systems and methods that create, provide, and facilitate new and interesting online auctions that are fair to each customer.\nSome online auction formats, such as the penny auction, also known as a click to bid auctions have been known to be subject to fraud, cheating, and abuse by both customers and auction operators. The two most common forms of cheating are known as “shill” bidders and “bot” bidders. A “shill” bidder is used by an auction operator to artificially inflate the price of the auction while increasing the amount of the customer bids by having employees of the auction operator pose as bidders and bid against the real customers. A “bot” bidder accomplishes the same result, however the false bidding is accomplished by software code, or script, which poses as a real person competitively bidding in the auction. The reputation of the penny auction industry has been severely damaged by the use of shill and bot bidders. The reputation of fraud in the industry has been so damaging that many people refuse to participate in a penny auction for fear of being cheated out of their money.\nAnother detractor to acquiring new customers and retaining existing penny auction customers is the potential for any customer to be out bid by a well-financed competitor customer, thereby winning most auctions. For example, a customer who has the personal financial ability to purchase ten times the bids of a lesser well-financed customer is at a distinct and nearly insurmountable advantage. The better financed customer could simply keep bidding in an auction against the lesser financed customer until he or she runs out of bids and is unable to purchase more bids to remain competing in the auction. This appearance that penny auctions unfairly provide advantage to well-funded customers is known to prevent many new potential customers from trying the product for a first time, thereby severely limiting the potential growth of a penny auction business.\nThe present invention is aimed at one or more of the problems identified above."}
{"text": "The present invention concerns a process and a device for validating/invalidating a message sent over an information transmission network.\nThe present invention is more particularly concerned with a process and a device for validating a message sent over a network between at least two stations connected to the network by means of a response in a communication frame, the message including at least one header field sent by a first station and one data field sent by a second station.\nIn prior art transmission processes and devices of the above kind, also known as communication frame response processes and devices, at least one of the stations accesses the network in order to transmit the beginning of a message frame, also known as the header field, and in accordance with the message header field another station accesses the network to transmit the end of the message frame, containing event or data information from that station, the end of the frame being also generally known as the data field.\nHowever, the above processes and devices give rise to a number of problems, in particular concerning updating and transmission of the data field. The time of transmission of the data is not under the control of the station producing the data, but rather that requiring the data.\nIf the station producing the data wishes to update it, there is a risk that the update will take place during the transmission of the data over the network, if the update is carried out in the meantime.\nThe risk is then that the corresponding station transmits internally inconsistent information, some of which has been updated and the remainder of which has not.\nOne prior art proposal for solving this problem is to add a validity marker to the data field, for example at the beginning of this field.\nThe marker is set to a particular value before any updating of the data and changed to a different value when the update has been done.\nA station reading the data can then analyze the marker to determine if the data is being modified and updated or not and therefore if it should be regarded as valid or invalid.\nHowever, this mechanism is inadequate as it does not solve the problem of transmitting an inconstant data field, in particular if the updates starts while the data is being transmitted.\nThe aim of the invention is therefore to solve these problems."}
{"text": "Conventional Petrie dishes used for procedures such as observation, growth, expansion, biopsy and manipulation of mammalian cells are not unique and are mostly generic. The user usually has to add droplets or media pools to the flat surface of the dish and then use it for the particular purpose. The media pools are thus randomly placed in the Petrie dish, and thus the technician must spend additional time locating the various media pools and observing the specimens in the individual media pools. The time spent can be minutes to tens of minutes. The standard Petrie dish thus complicates the process of monitoring specimen growth and development, and complicates the ability to tell which specimen is which.\nOne problem with using generic Petri dishes in specimen procedures is that they do not have unique dish configurations which enhance their utility in such procedures. Another problem with the current dishes is that they require the user to use micro drops of a specimen medium solution on the surface of the dish. Some dishes may create surface tension problems with the specimen medium solution. This can result in the micro drops collapsing and the media solution that the sample is in then becomes compromised by an overlain layer of oil. The reverse may happen too and the droplets may not adhere well to the bottom of the dish causing them to move in the dish.\nIt would be highly desirable to provide a specimen treatment dish assembly that will inherently locate the specimen being treated in a definite predetermined location in the dish and retain the specimen in that definite location. It would be further desirable to provide a specimen treatment dish assembly which does not require movement of the optics, such as a microscope, to observe and log the changes in the several specimens being treated in the dish."}
{"text": "Research in the area of flapping wing flight has been conducted by the biological and engineering communities. Researchers have studied the flight dynamics of animals and ornithopters. However, none of these efforts has resulted in an aircraft capable of six-degrees-of-freedom flight with only two actuators."}
{"text": "The present invention relates to compositions useful for inducing alignment of a liquid crystal medium with polarized light and liquid crystal display elements.\nMost liquid crystal devices, including displays, have a finite pre-tilt angle, controlled, for instance, by the mechanical buffing of selected polymeric alignment layers. Liquid crystal molecules in contact with such a layer align parallel to the buffing direction, but not exactly parallel to the substrate. The liquid crystal molecules are slightly tilted from the substrate, for instance, by about 0.5-15 degrees. For most display applications, a finite and uniform pre-tilt angle of the liquid crystal is desirable. For other display applications requiring predominately homeotropic alignment of liquid crystals, a much higher pre-tilt angle is desired, for instance, 85-89.9 degrees. In these cases, controlling the pre-tilt by a mechanical rubbing process is very difficult.\nA new process for aligning liquid crystals with polarized light was disclosed in U. S. Pat. No. 4,974,941 entitled xe2x80x9cProcess of Aligning and Realigning Liquid Crystal Mediaxe2x80x9d. Controlling the pre-tilt angle of liquid crystals in contact with optical alignment layers, while maintaining high uniformity of alignment, is an important requirement for optical alignment materials. Continuing effort has been directed to the development of compositions for optical alignment of liquid crystals and liquid crystal displays.\nThe present invention provides polarizable fluorinated amines of the general formula \nwherein A is a divalent or trivalent organic moiety, Lf is a divalent organic radical comprising at least four fluorinated carbon atoms; P is a polar group comprising a Pi electron system containing at least one heteroatom selected from the group N, O, and S; and s is 1-2.\nThe invention further provides polyamic acids prepared from these amines and at least one tetracarboxylic dianhydride and the corresponding polyimides derived therefrom.\nThe invention further provides a process for inducing pre-tilt in alignment of a liquid crystal using an optical alignment layer bonded to such polyimides, and liquid crystal display elements prepared from such optical alignment layers."}
{"text": "The invention relates generally to liquid phase liquid crystal polymers having imino- and amino-borane backbone chains with side chains that impart desirable properties and applications."}
{"text": "The present invention relates to a bumper system for motor vehicles.\nAs is known, the main function of bumper systems for motor vehicles is to convert impact energy to deformation energy in a collision with another vehicle or with any stationary or non-stationary obstacle. This limits damage to the motor vehicle.\nAs a rule, one known embodiment of such a bumper system comprises a cross-member that has a U-shaped cross-section, with the bar portion thereof connected via crush boxes to the longitudinal supports of the motor vehicle. The ends of the leg portions of the cross-member are connected, especially joined, using a closure plate, such that the closure plate prevents the cross-member from opening and/or flattening under a bending load. The profile of the bumper system is closed and thereby achieves high flexural rigidity, while remaining lightweight.\nBoth steel alloys and aluminum alloys are suitable materials for the cross-member and the closure plate. Furthermore, combinations of these materials are possible."}
{"text": "1. Field of the Invention\nThe invention relates to a device for the region of cans held under propellant pressure where liquid to pasty substances containing solvent pass out, including an operating device for delivering the substances.\nIt is known to store substances initially in cans, in which they are held under pressure. This pressure can be produced in various ways, for example by means of a propellant gas, bladders under internal pressure, pistons under spring pressure, or the like. When an operating device is operated, the substance passes out into the atmosphere. For this purpose, a spring-loaded button is pressed. A part of this actuating device can also be tilted by pulling a lever 36, so that the substance passes out of the tube 37, as described in European Patent Application 91 114 634.8. The substances may, for example, be lacquers, lubricants and, in particular, substances which in the air harden to a foam.\n2. Discussion of Prior Art\nIt has become possible, in accordance with Swiss Patent Document CH-PS 676 354, to produce the valve part of the operating device without rubber or other comparable substances. Such devices can consequently be stored and transported even in a horizontal position. However, the weakest link in the chain is still the region downstream of the valve. The substance, particularly the one which subsequently forms the polyurethane foam, can there form a clot. Although the problems in the valve region have been overcome, it is therefore possible to use the can only once if the substance contained in the so-called stem loses solvents through evaporation. In addition, moisture can pass from the interior to the exterior of the device, and also in particular from the exterior to the interior.\nParticularly in the field of hobbies it is known that small amounts of substances are needed. If larger amounts are bought, this is detrimental in various respects: more than is needed is bought. The unused residues must be disposed of as special refuse. To some extent this is also the case when smaller cans are produced. It would, however, be more advantageous also to make larger cans available for hobby purposes, provided they can be used until they are quite empty. It would be ideal if, over a long period of time, it were possible to use the can at the end of its life in exactly the same way as when it is fresh from the factory.\nIf the substance has already solidified in the stem, it is also possible that when the valve is subsequently operated a number of times, solidified material will also pass into the valve, so that the latter itself will leak. Even if the removable part of the operating device were to be cleaned, the leaking valve could in the meantime have entirely or to a substantial extent destroyed the pressure in the can."}
{"text": "Without limiting the scope of the invention, its background is described in connection with the protection of integrated circuit fingerprint sensors from the environment during regular use, as an example.\nHeretofore, in this field, the detection of fingerprint patterns has been useful for the identification of specific individuals based on the observation that each individual person has a unique fingerprint. Fingerprints, therefore, can be used not only to positively identify individuals, but to exclude individuals whose fingerprint profile does not match a pre-existing set of patterns.\nFingerprint sensing has evolved from optical and mechanical sensing technology that acquires a fingerprint image. Generally, the mechanical and optical sensors obtained a fingerprint image using a scanner or a camera, processed the acquired information into an analog or digital signal that could be analyzed, and provided an output based on the acquired signal. Unfortunately, the lighting and contrast conditions available at the time the image was acquired affected the analysis and consequently the output from the sensor.\nAnother class of fingerprint sensors are capacitive sensors, such as that disclosed in U.S. Pat. No. 4,353,056 issued to Tsikos. The Tsikos patent demonstrates the use of a sensor that incorporates a sensing member that has a sensing surface for receiving a fingerprint. The sensing surface has a means for sensing the ridges and valleys of the skin of the finger under observation. A sensing member containing a multitude of capacitors that sense the patterns of the fingerprint when the finger is pressed against the sensing surface. The information obtained is transformed into an electric signal. The capacitors are insulated from the environment of use by a flexible membrane that conforms itself to the contour of the fingerprint. Unfortunately, the repeated cycles of flexing and compression of the flexible membrane can lead to device failure and the need to replace the membrane.\nU.S. Pat. No. 4,353,056 issued to Ruell, et al., discloses a fingerprint sensor that provides an electric output signal in response to the topographic relief of the fingerprint. The sensor incorporates a contact body that is formed at least in part by a light transparent elastic material. The elastic contact material can be attached to a flat sensor plate that has a light receiving surface. The sensor also incorporates a light source and a photodetector to measure the valleys and ridges of the fingerprint. As with the sensor described hereinabove, the elastic nature of the contact body leads to cycles of compression and flexing that lead to the deterioration of the contact point between the sensor and the finger."}
{"text": "1. Field of the Invention\nThe present invention relates to an air intake system of an aircraft and having an excellent secretiveness (stealth).\n2. Description of the Prior Art\nHeretofore it has been known that if an air intake port of an air intake system of an aircraft is provided on an upper surface of an airframe of the aircraft, the secretiveness or security of the aircraft will be high, especially against radair beam irradiation from below, i.e., on the lower surface of the airframe. However, since turbulence of air on the upper surface of the airframe is violent and the performance at such an air intake port is degraded upon taking a large angle of attach and at high mobility, it was difficult to practically employ such system.\nIn other words, in the prior art if it is contemplated to dispose an air intake port having excellent secretiveness on an upper surface of an airframe, when the aircraft takes a large angle of attack and at high mobility, the air intake port would introduce to the engine tubulent air on the upper surface of the airframe, and thus the performance (total pressure recovery, distortion, etc.) of the air intake port would be degraded. Under such circumstances, the engine would stop frequently, resulting in fatal shortcomings for an airframe having high mobility."}
{"text": "Air filtration masks (referred to herein as “filter masks”) are widely used to protect people from air borne contaminants and gasses. For example, air borne dust particles are a known hazard commonly on work sites. Consequently, workers normally wear filter masks to avoid inhaling the dust particles. To that end, filter masks used in this application are manufactured with a filter material specified to prevent, among other things, a substantial majority of dust particles from being inhaled by the worker.\nIn addition to primarily filtering inhaled air, some filter masks are specifically manufactured to filter both inhaled and exhaled air. For example, hospital staff often wear filter masks to prevent both their germs from infecting patients, and patients' germs from infecting them.\nThere is a need in the art to improve the filtration efficiency of filter masks. Accordingly, filter masks with a higher efficiency filter layer and/or multiple filter layers have been developed for that purpose. However, this often has the undesirable effect of increasing the air resistance through the filter mask and may cause several problems.\nFor example, a person wearing the filter mask may have a more difficult time breathing due to the increased air resistance. To overcome this problem while still providing improved filtration efficiency, filter masks have been developed that have an increased filter area. Manufacture of such filter masks, however, can be quite complex. For example, increasing the filter area can cause various portions of the filter layer to overlap or can be costly to construct. Overlap can effectively increase the thickness of the filter layer, thus causing the same air resistance problem as discussed above.\nAdditionally, since a person wearing the mask while performing manual labor must typically breathe heavier, the filter layer(s) is more likely to flex and eventually collapse around the face. This collapse may cause portions of the face mask to contact and irritate the face of the person wearing the face mask, as well as cause discomfort. Consequently, efforts have been made to stiffen the mask, such as by adding additional material to the filter mask. However, adding additional material to the face mask adds complexity to the production process and increases cost."}
{"text": "Field of the Invention\nThe present invention generally relates to a vehicle door structure. More specifically, the present invention relates to vehicle door structure having a first bracket and a reinforcing bar with a contoured end surface that is spaced apart from adjacent surfaces of the first bracket in an undeformed state, the contoured end surface being shaped to maximize surface area contact between the first end section of the reinforcing bar and the adjacent surfaces of the first bracket in response to a force being applied to the reinforcing bar and causing the first bracket to deform.\nBackground Information\nMany areas of a vehicle body structure are provided with reinforcing structures. For example, vehicle doors usually include reinforcing structures."}
{"text": "Atrial tachyarrhythmias are quite common, and are seen even in patients who are not acutely ill. One such atrial tachyarrhythmia, atrial fibrillation, does not usually pose an immediate threat to the patient's life, but prolonged atrial fibrillation can compromise patient hemodynamics. Inadequate blood flow due to atrial fibrillation may lead to blood clots, which increase the risk of strokes. Other symptoms may include palpitations of the heart, dizziness, or loss of consciousness.\nOne method of treating atrial fibrillation is to deliver a countershock to terminate the arrhythmia. For example, the countershock may be applied via an external or implantable defibrillator. However, this method of treating atrial fibrillation risks causing ventricular tachyarrhythmias, including life-threatening ventricular fibrillation.\nOne method of reducing the risk of ventricular fibrillation is by synchronizing the delivery of a countershock to an electrical ventricular activation (R-wave) of the heart. There is, however, a need in the art to further reduce the risk of inducing ventricular tachyarrhythmias, such as life-threatening ventricular fibrillation, during the countershock treatment of comparatively less exigent atrial tachyarrhythmias."}
{"text": "1. Field of the Invention\nThis invention relates to a method for preparing granules having a uniform granule size, good disintegrating properties, relatively great apparent density and abrasion resistance by subjecting a crystalline sugar, a basis, adjuvants and a binder-containing solution to a mixing apparatus, whereby the crystalline sugar is coated with the basis and adjuvants to form granules of a uniform particle size.\n2. Prior Art\nIn general, granule or particle products have been heretofore prepared by powdering raw materials such as a basis and some suitable adjuvants, mixing the powdered material with a corrigent such as powdered sugar under wetting conditions, and granulating the mixture by means of an extruder-type or crushing type granulator.\nThe just-mentioned prior-art method is, however, defective in that: (1) The yield of the granule product becomes low due to substantial loss of the powdered material in the binding stage and, if a binder is used in a larger amount so as to raise the yield, the ultimate granular product will be deteriorated in disintegratability and solubility; (2) It is difficult to produce granules of uniform size and relatively great apparent density, and the commercial value of the products is lowered due to undesirably wide distribution of granule sizes and inconveniences in handling; and (3) A great number of production steps and manual operations are required, causing frequent operational troubles.\nIn order to overcome these defects, many improved methods have been proposed by introducing new apparatus of different types, including the method using a fluidized bed granulation apparatus and the method of coating core material with a sprayed suspension with use of a coating apparatus. However, the last-metioned coating method has a vital disadvantage that it requires a solvent in such a large amount that may spoil the working conditions with increased possibility of risky accidents, such as fire, explosion and the like, and also may give rise to environmental pollution. In addition, an enormous cost is required for installation of the facilities for collecting the solvent in order to prevent environmental pollution. Moreover, the spraying method has also the disadvantage that it is difficult to coat the core material completely with a powdery material. This results in reduction in yield and undesirably requires provision of a fine powder-collecting apparatus. Furthermore, the coated material is easily excoriated from the ultimate product even by slight friction, such as, for example, that occurred by sieving, due to insufficient binding between the core material and the coat, thus spoiling the commercial value of the product."}
{"text": "1. Field of the Present Invention\nThe present invention relates generally to database systems, and more particularly to systems and methods for the identification of latent relationships amongst data elements in very large computer databases.\n2. Description of Related Art\nBusinesses, governments, agencies, and other institutions have for a long-time required various people to fill out reports and forms. The fate of many of these filings is they get forgotten. Paper forms are, in general, hard to file and access. So very often the purpose and objective in requiring the filing of forms gets lost in the impracticality of dealing with so much paper.\nThe Federal suspicious activity report (SAR) is a case-in-point. A SAR is required to be filed by bank officials in every instance in which some unusual financial transaction or contact has occurred. Under Title 12 Code of Federal Regulations xc2xa721 (12 CFR 21), all financial institutions operating in the United States, including insured banks, savings associations, savings association service corporations, credit unions, bank holding companies, non-bank subsidiaries of bank holding companies, Edge and Agreement corporations, and American branches and agencies of foreign banks, are required to make this report following the discovery of insider abuse involving any amount, violations aggregating $5,000 or more where a suspect can be identified, violations aggregating $25,000 or more regardless of a potential suspect, or transactions aggregating $5,000 or more that involve potential money laundering or violations of the Bank Secrecy Act.\nThe information provided by SAR filings across the country is invaluable to law enforcement. But the data entry, indexing and filing of paper reports results in a lot of information being lost due to misfilings, alternative wordings, misspellings, typographical errors, and other inconsistencies. So the Treasury Department, has a website (www.occ.treas.gov) that provides downloadable software for electronic filing and distribution of SAR\"\"s. To make that report, the filing institution prepares a SAR, and files it with the Financial Crimes Enforcement Network (FinCEN) of the Department of the Treasury through the IRS Detroit Computing Center. The reports are then made available electronically to appropriate law enforcement agencies.\nBanks that do not wish to file electronically can continue to file paper-based reports. FinCEN has made copies of the forms available for download in Adobe Acrobat portable document format (PDF).\nThe FinCEN SAR database and other modern databases are proliferating in almost every area of science, technology, and society in general. Information systems are getting larger, more complex, and more important. Data mining technology has been developed that extracts useful information from databases. The identification of latent relationships among data elements in very large databases is a particularly lofty goal of current systems. This problem is made even more difficult when there are errors or variations in the type of data such as misspellings and abbreviations.\nA variety of data mining methods have appeared on the commercial market in database software packages. But these prior art products are only useful when applied to small databases, where the data is pristine, and when identifying the sought-for information is comparatively straightforward. Such conventional database systems are confounded as the size of the database, the noise in the data, and the complexity of the information increases.\nHigh-performance data mining tools take advantage of the format and structure of a given set of data. The data format and data structure are usually not subject to change. Data users must understand the data well enough to be able to identify any relationships that may exist amongst the data elements, even though the nature of such relationships may not be known beforehand. Search speeds can be improved if each database is decomposed into separate components. Once decomposed, the detection of latent relationships amongst the seemingly independent data elements becomes easier, and the speed at which the database may be searched for information can be significantly enhanced.\nMany databases use fixed-width data fields, while others use a combination of fixed fields, free text and even graphic images. Fixed-field databases are the easiest to deal with, and conventional approaches to data mining may produce reasonable results provided the data is error-free. Free-text databases are more troublesome, particularly when any latent relationship information occurs in variable-width fields. Misspelled words and the use of abbreviations only complicates the job more.\nAn object of the present invention is to provide a database method and system for the identification of any latent relationships amongst independent data elements in very large databases.\nA further object of the present invention is to provide a database method and system for effective data mining by relatively inexperienced users.\nAnother object of the present invention is to provide a law enforcement tool for detecting obscure relationships and subtle activities of criminals and their organizations.\nBriefly, a database method embodiment of the present invention quickly identifies latent relationships among data elements in very large databases. The database method converts each word in a document in a large database into a very wide digital-bit vector. Unique words appearing in the document are encoded into particular bits of the corresponding vector by writing ones. The possibility of a keyword being present in any document can therefore be quickly ascertained by throwing out all vectors with digital zeros in the critical bit positions that must be ones.\nAn advantage of the present invention is that a database search tool is provided that can scan exceedingly large databases quickly and efficiently.\nAnother advantage of the present invention is a data mining tool is provided that can cope with alternate spellings, misspellings, and abbreviations used for keywords in report form inputs.\nThese and many other objects and advantages of the present invention will no doubt become obvious to those of ordinary skill in the art after having read the following detailed description of the preferred embodiments which are illustrated in the various drawing figures."}
{"text": "1. Field of the Invention\nThe present invention generally relates to remotely operated systems, and more particularly to a computerized system for monitoring, reporting on, and controlling remote systems by transferring information signals through a wide area network (WAN) and using software applications hosted on a connected server to appropriately process the information.\n2. Discussion of the Related Art\nAs is known, there are a variety of systems for monitoring and controlling manufacturing processes, inventory systems, emergency control systems, and the like. Most automatic systems use remote sensors and controllers to monitor and automatically respond to system parameters to reach desired results. A number of control systems utilize computers to process system inputs, model system responses, and control actuators to implement process corrections within the system. Both the electric power generation and metallurgical processing industries have had success controlling production processes by implementing computer controlled control systems in individual plants.\nOne way to classify control systems is by the timing involved between subsequent monitoring occurrences. Monitoring processes can be classified as aperiodic or random, periodic, and real-time. A number of remotely distributed service industries implement the monitoring and controlling process steps through manual inspection and intervention.\nAperiodic monitoring systems (those that do not operate on a predetermined cycle) are inherently inefficient as they require a service technician to physically traverse an area to record data, repair out of order equipment, add inventory to a vending machine, and the like. Such service trips are carried out in a number of industries with the associated costs being transferred to the consumers of the service.\nConversely, utility meter monitoring, recording, and client billing are representative of a periodic monitoring system. In the past, utility providers sent a technician from meter to meter on a periodic basis to verify meter operation and to record utility use. One method of cutting operating expenses in the utility industry involved increasing the period at which manual monitoring and meter data recording was performed. While this method decreased the monitoring and recording expense associated with more frequent meter observation and was convenient for consumers who favor the consistent billed amounts associated with “budget billing,” the utility provider retained the costs associated with less frequent meter readings and the processing costs associated with reconciling consumer accounts.\nLastly, a number of environmental and safety systems require constant or real-time monitoring. Heating, ventilation, and air-conditioning systems, fire reporting and damage control systems, alarm systems, and access control systems are representative systems that utilize real-time monitoring and often require immediate feedback and control. These real-time systems have been the target of control systems theory and application thereof for some time.\nA problem with expanding the use of control systems technology to distributed systems are the costs associated with the sensor-actuator infrastructure required to monitor and control functions within such systems. The typical approach to implementing control system technology is to install a local network of hard-wired sensors and actuators along with a local controller. Not only is there expense associated with developing and installing appropriate sensors and actuators but the added expense of connecting functional sensors and controllers with the local controller. Another prohibitive cost associated with applying control systems technology to distributed systems is the installation and operational expense associated with the local controller.\nAccordingly, an alternative solution to applying monitoring and control system solutions to distributed systems that overcomes the shortcomings of the prior art is desired."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a light source and, in particular, to a light-emitting diode (LED) package.\n2. Description of Related Art\nAs the luminous efficiency of light emitting diodes (LEDs) is increasingly improved, LEDs have replaced fluorescent lamps and incandescent lamps in some fields, for example, lamps of scanners requiring for quick response, lamps of projection apparatus, backlight sources or front light sources of liquid crystal displays (LCDs), illumination for dashboards of automobiles, traffic lights, common illumination devices, etc. Compared with conventional lamps, the LEDs have absolute advantages, for example, small volume, long lifespan, low driving voltage/current, non-fragile, mercury free (no pollution), and good luminous efficiency (power saving).\nFIG. 1 is a schematic structural view of a conventional LED package. Referring to FIG. 1, the conventional LED package 100 includes an LED chip 110, a lead frame 120, and an encapsulant 130. The LED chip 110 is disposed on the lead frame 120 and electrically connected therewith. The encapsulant 130 wraps the LED chip 110 and a part of the lead frame 120 for protecting the LED chip 110. The encapsulant 130 has a convex top surface 132 for preventing the total internal reflection (TIR) of the beam 112 emitted from the LED chip 110 due to the interface between the encapsulant 130 and air.\nFIG. 2 is a schematic structural view of another conventional LED package. Referring to FIG. 2, the other conventional LED package 200 includes an LED chip 210, a circuit board 220, and an encapsulant 230. The LED chip 210 is disposed on the circuit board 220 and electrically connected therewith. The encapsulant 230 is disposed on the circuit board 220 and wraps the LED chip 210. The encapsulant 230 also has a convex top surface for preventing the TIR of the beam 212 emitted from the LED chip 210.\nHowever, referring to FIGS. 1 and 2, the range of the light emitting angle θ1 and θ2 of the conventional LED packages 100 and 200 is large, approximate to 180 degrees, such that the LED packages 100 and 200 are not suitable for illuminating distant objects. For example, when the LED package 100 or 200 is applied as a light source in the projection apparatus (not shown), the beam 112 or 212 is hard to be mostly converged on the light valve (not shown) of the projection apparatus, such that the brightness of the light valve is not enough."}
{"text": "The present invention relates to personal armor for protection against projectile threats in general, and more particularly to armor intended to counter higher speed projectiles such as rifle rounds.\nPersons exposed to projectile threats, such as police officers and soldiers, may seek a certain level of protection by wearing armored clothing. Such body armor is available in a variety of forms to address various levels of potential threats.\nLow velocity projectiles such as handgun rounds, fragmentation rounds from a grenade or mortar, and miscellaneous shrapnel may be countered by so-called “soft armor.” Soft armor is worn in the form of jackets, vests, etc. which are composed of assemblies of ballistic fabric such as those formed from DuPont Kevlar® fibers. Typically, soft armor will employ 15-30 layers of thin flexible tightly woven ballistic fabric. This material is generally lightweight and effective in stopping handgun rounds. However, the multiple layers required to make the soft armor are also an effective insulator, trapping heat close to the wearer's body.\nIn a more serious threat situations, where higher velocity rifle rounds must be countered, soil armor has typically been supplemented with hard armor. The hard armor is fabricated of rigid plates of ceramic, polymer, or metal. These plates are usually molded or formed to conform to a standard wearer. To provide the greatest area of coverage, and hence protection, these plates are often large and hence can interfere with the wearer's mobility. A common approach to mounting the plates to the wearer is to secure them within exterior pockets fabricated on a soft armor jacket or vest. This modular approach allows the wearer to assess the level of risk and to add or subtract hard armor as the situation demands. However, the result is a more bulky and potentially hot armor assembly. Because the hard armor plates can compromise mobility, the wearer may choose to sacrifice a greater level of protection for greater ease of movement. The option to remove some levels of protection, however, means that occasionally the wearer of the armor encounters a higher velocity projectile without the appropriate hard armor protection.\nWhat is needed is a body armor system which is effective against high velocity projectiles yet which does not unduly compromise wearer mobility or temperature regulation."}
{"text": "Mobile communication systems have been developed to provide voice services while ensuring the activity of a user. However, the mobile communication systems have been expanded to their regions up to data services as well as voice. Today, the shortage of resources is caused due to an explosive increase of traffic, and more advanced mobile communication systems are required due to user's need for higher speed services.\nRequirements for a next-generation mobile communication system basically include the acceptance of explosive data traffic, a significant increase of a transfer rate per user, the acceptance of the number of significantly increased connection devices, very low end-to-end latency, and high energy efficiency. To this end, research is carried out on various technologies, such as dual connectivity, massive Multiple Input Multiple Output (MIMO), in-band full duplex, Non-Orthogonal Multiple Access (NOMA), the support of a super wideband, and device networking."}
{"text": "This invention relates to a reagent and method for determining the activity of retroviral protease, e.g. human immunodeficiency virus (HIV) protease. More particularly, the invention relates to novel fluorogenic substrates and a fluorometric assay for HIV protease and other retroviral proteases.\nAcquired immune deficiency syndrome, which only a few years ago was a medical curiosity, is now a serious disease. As a consequence, a great effort is being made to develop drugs and vaccines to combat AIDS. The AIDS virus, first identified in 1983, has been described by several names. It is the third known T-lymphocyte virus (HTLV-III) and has the capacity to replicate within cells of the immune system and thereby lead to a profound destruction of T4.sup.+ T-cells (or CD4.sup.30 cells). See, e.g., Gallo et al., Science 224, 500-503 (1984), and Popovic et al., Ibid., 497-500 (1984). This retrovirus has been known as lymphadenopathy-associated virus (LAV) or AIDS-related virus (ARV) and, most recently, as human immunodeficiency virus (HIV). Two distinct AIDS viruses, HIV-1 and HIV-2, have been described. HIV-1 is the virus originally identified in 1983 by Montagnier and co-workers at the Pasteur Institute in Paris [Ann. Virol. Inst. Pasteur 135 E, 119-134 (1984)], while HIV-2 was more recently isolated by Montagnier and his co-workers in 1986 [Nature 326, 662 (1987)]. As used herein HIV is meant to refer to these viruses in a generic sense.\nA useful approach being investigated recently for potential use in the treatment of AIDS is the development of synthetic peptides as inhibitors of retroviral protease. Thus, it is known that retroviruses, including the human immunodeficiency virus (HIV), express their genetic content by directing the synthesis of a polyprotein by the host. This precursor is then processed by proteolysis to give essential viral enzymes and structural proteins. A virally encoded enzyme, the retroviral protease, is contained within the polyprotein and is responsible for the specific cleavages of the polyprotein yielding mature viral proteins.\nVarious novel peptide inhibitors of retroviral protease such as HIV protease are described in applicants' copending application Ser. No. 07/320,742, filed Mar. 8, 1989. Included among these are short peptides of about 4 to 8 amino acid residues derived from HIV cleavage sites such as, e.g. p24/p15, which have been modified to incorporate an internal CH.sub.2 NH bond isostere. An illustrative example of such peptides is\nAc--Thr--Ile--Nle--.psi.[CH.sub.2 NH]--Nle--Gln--Arg--NH.sub.2.\nHIV protease thus is a logical therapeutic target in the search for AIDS drugs. In order to facilitate the exploration of potential inhibitors of HIV protease and to optomize lead compounds, a convenient assay to kinetically characterize the inhibitor which allows for automation and high throughput would be desirable. Initial assays used in the characterization of HIV protease were based on HPLC separation of products and substrate at fixed time intervals. See Moore et al., Biochem. Biophys. Res. Commun. 159, 420-425 (1989); Darke et al., Ibid. 156, 297-303 (1988); Schneider and Kent, Cell 54, 363-368 (1988). The HPLC assay quickly became the rate-limiting step in the development of HIV protease inhibitors with desirable therapeutic properties. A continuous assay which allows for quantitative kinetic characterization of the interaction of the inhibitors with HIV protease would be preferable. The strategy of developing either a chromogenic or fluorogenic substrate is proposed by the present inventors based on their work characterizing short peptide substrates of known HIV protease cleavage sites. Toth et al., Proc. 11th Am. Peptide Symp., Rivier and Marshall, Eds., ESCOM, Leiden, In Press 1990. Fluorogenic substrates for hydrolytic enzymes have been widely used in biochemistry because of their high sensitivity. See, e.g., Guilbault, Enzymatic Methods of Analysis, Pergamon Press, 1970, pp. 43-47.\nRecently, Nashed et al., Biochem. Biophys. Res. Commun. 163, 1079-1085 (1989), and Hyland et al., Abstr. 2nd Int. Conf. Drugs Res. Immunologic and Infectious Disease. AIDS. New York Acad. Sci., I-20 (1989), have reported a spectrophotometric assay based on cleaved chromogenic substrates containing p--NO.sub.2 --Phe at position P1. In these cases, the substances were an octapeptide, a nonapeptide, and a decapeptide. A high-throughput, radiometric assay for screening for HIV protease inhibitors, which is not continuous, has also been reported by Hyland et al., supra."}
{"text": "The present apparatus is directed to a fluid mud pump and, more particularly, to a mud pump to be utilized to intensify fluid pressure for use in drilling oil wells or in conditioning oil wells such as fracturing with extremely high pressure or abrasive fluids. Various mud pumps and pressure intensification pumps are already known to exist that employ various and sundry means to overcome the difficulties encountered in prolonged pumping of high volume, high pressure, and abrasive materials. The present invention is an apparatus which will provide improvement in mud pumping operations in such areas as reduced mud pressure pulsation, less operating energy required for fluid pressure intensification, slower operating piston speeds and longer piston strokes thus resulting in extended life of all operating parts, wider range of mud flow and pressure controllability, greater simplicity of manufacture, improved adaptability and operation, plus other less apparent improvements. Thus the context of the problem to be dealt with in the present invention is that of a non pulsating output, highly efficient and controllable hydraulic powered fluid output pump."}
{"text": "Many different types and styles of memory exist to store data for computers and similar type systems. For example, random access memory (RAM), dynamic random access memory (DRAM), static random access memory (SRAM), read only memory (ROM), programmable read only memory (PROM), electrically programmable read only memory (EPROM), electrically erasable programmable read only memory (EEPROM) and flash memory are all presently available to accommodate data storage.\nEach type of memory has its own particular advantages and disadvantages. For example, DRAM and SRAM allow individual bits of data to be erased one at a time, but such memory loses its data when power is removed. EEPROM can alternatively be easily erased without extra exterior equipment, but has reduced data storage density, lower speed, and higher cost. EPROM, in contrast, is less expensive and has greater density but lacks erasability.\nFlash memory, has become a popular type of memory because it combines the advantages of the high density and low cost of EPROM with the electrical erasability of EEPROM. Flash memory can be rewritten and can hold its contents without power, and thus is nonvolatile. It is used in many portable electronic products, such as cell phones, portable computers, voice recorders, etc. as well as in many larger electronic systems, such as cars, planes, industrial control systems, etc. Flash memory is generally constructed of many memory cells where single bits of data are stored in and read from respective memory cells.\nIndividual memory cells are generally organized into individually addressable units or groups, which are accessed for read, program, or erase operations through address decoding circuitry. The individual memory cells are typically comprised of a semiconductor structure adapted for storing a bit of data. For instance, many conventional memory cells include a stacked gate metal oxide semiconductor (MOS) device, such as a transistor in which a binary piece of information may be retained. The memory device includes appropriate decoding and group selection circuitry, as well as circuitry to provide voltages to the cells being operated on.\nThe erase, program, and read operations are commonly performed by application of appropriate voltages to certain terminals of the memory cell. In an erase or write operation the voltages are applied so as to cause a charge to be removed or stored in the memory cell. In a read operation, appropriate voltages are applied so as to cause a current to flow in the cell, wherein the amount of such current is indicative of the value of the data stored in the cell. The memory device includes appropriate circuitry to sense the resulting cell current in order to determine the data stored therein, which is then provided to data bus terminals of the device for access by other devices in a system in which the memory device is employed.\nThe traditional memory cell generally has a source, a drain, and a channel in a substrate or P-well, as well as a stacked gate structure overlying the channel. The stacked gate may further include a thin gate dielectric layer (sometimes referred to as a tunnel oxide) formed on the surface of the P-well. The stacked gate also includes a polysilicon floating gate overlying the tunnel oxide and an interpoly dielectric layer overlying the floating gate. The interpoly dielectric layer is often a multilayer insulator such as an oxide-nitride-oxide (ONO) layer having two oxide layers sandwiching a nitride layer. Lastly, a polysilicon control gate overlies the interpoly dielectric layer.\nIn a NOR configuration, the control gate is connected to a wordline associated with a row of memory cells to form sectors of such cells. In addition, the drain regions of the cells are connected together by a conductive bitline. The channel of the cell conducts current between the source and the drain in accordance with an electric field developed in the channel by the stacked gate structure. Respective drain terminals of the transistors within a single column are connected to the same bitline. In addition, respective flash cells associated with a given bitline have stacked gate terminals coupled to a different wordline, while all the flash cells in the array generally have their source terminals coupled to a common source terminal. In operation, individual flash cells are addressed via the respective bitline and wordline using the peripheral decoder and control circuitry for programming (writing), reading or erasing functions.\nBy way of further detail, the single bit stacked gate flash memory cell is programmed by applying a relatively high voltage to the control gate and connecting the source to ground and the drain to a predetermined potential above the source. A resulting high electric field across the tunnel oxide leads to the Fowler-Nordheim tunneling phenomena. During this process, electrons in the core cell channel region tunnel through the gate oxide into the floating gate and become trapped in the floating gate since the floating gate is surrounded by the interpoly dielectric and the tunnel oxide. As a result of the trapped electrons, the threshold voltage of the cell increases. This change in the threshold voltage (and thereby the channel conductance) of the cell created by the trapped electrons is what causes the cell to be programmed.\nIn order to erase a typical single bit stacked gate flash memory cell, a relatively high voltage is applied to the source, and the control gate is held at a negative potential, while the drain is allowed to float. Under these conditions, a strong electric field is developed across the tunnel oxide between the floating gate and the source. The electrons that are trapped in the floating gate flow toward and cluster at the portion of the floating gate overlying the source region and are extracted from the floating gate and into the source region by way of Fowler-Nordheim tunneling through the tunnel oxide. As the electrons are removed from the floating gate, the cell is erased.\nFor a read operation, a certain voltage bias is applied across the drain to source of the cell transistor. The drain of the cell is the bitline, which may be connected to the drains of other cells in a byte or word group. The voltage at the drain in conventional stacked gate memory cells is typically provided at between 0.5 and 1.0 volts in a read operation. A voltage is then applied to the gate (e.g., the wordline) of the memory cell transistor in order to cause a current to flow from the drain to source. The read operation gate voltage is typically applied at a level between a programmed threshold voltage (Vt) and an unprogrammed threshold voltage. The resulting current is measured, by which a determination is made as to the data value stored in the cell.\nA very modern memory technology is dual bit memory, which allows multiple bits to be stored in a single cell. In this technology, a memory cell is essentially split into two identical (mirrored) parts, each of which is formulated for storing one of two independent bits. Each dual bit memory cell, like a traditional cell, has a gate with a source and a drain. However, unlike a traditional stacked gate cell in which the source is always connected to an electrical source and the drain is always connected to an electrical drain, respective dual bit memory cells can have the connections of the source and drain reversed during operation to permit the storing of two bits.\nIn a virtual ground type architecture, dual bit memory cells have a semiconductor substrate with implanted conductive bitlines. A multilayer storage layer, referred to as a “charge-trapping dielectric layer”, is formed over the semiconductor substrate. The charge-trapping dielectric layer can generally be composed of three separate layers: a first insulating layer, a charge-trapping layer, and a second insulating layer. Wordlines are formed over the charge-trapping dielectric layer substantially perpendicular to the bitlines. Programming circuitry controls two bits per cell by applying a signal to the wordline, which acts as a control gate, and changing bitline connections such that one bit is stored by the source and drain being connected in one arrangement and a complementary bit is stored by the source and drain being connected in another arrangement.\nA continuing trend in the electronics industry is to scale down electronic devices to produce smaller, yet more powerful devices (e.g., cell phones, digital cameras, etc.) that can perform a greater number of increasingly complex functions faster and with less power. To achieve this, semiconductors and integrated circuits (e.g., memory cells, transistors, etc.) utilized in these devices are continually reduced in size. The ability to “pack” more of these components onto a single semiconductor substrate, or a portion thereof (known as a die), also improves fabrication efficiency and yields. Accordingly, there are ongoing efforts to, among other things, increase the number of memory cells that can be fabricated onto individual semiconductor wafers (or die).\nOne technique to pack more memory cells/transistors into a smaller area is to form their structures and component elements closer together. Forming bitlines closer together, for example, shortens the length of a channel defined there-between and allows more devices to be formed in the same area. This can, however, cause certain phenomena to become more prevalent and have a substantial impact on device performance. For example, Vt roll-off may become more pronounced as respective channel lengths are decreased. In particular, a threshold voltage (Vt), or voltage at which a transistor or memory cell “turns on” or becomes active to store and/or erase a bit of data (e.g., a charge), may vary (e.g., become reduced) as the respective channel length is decreased (e.g., source and drain regions are brought closer together). Additionally, the severity of Vt roll-off may vary among different memory cells. In this manner, not all of a selected group of memory cells may be programmed/erased when the same or a similar voltage is applied to the cells. This could result in corrupt data and/or inaccurately stored or programmed data.\nSimilarly, other issues, such as leakage currents, cross talk, etc. can also result from scaling. For example, isolating two bits or charges stored in a charge trapping layer becomes increasingly difficult as the channel length is decreased and the bits are brought closer together. In this manner, cross-talk can occur, the bits can contaminate one another and operations performed on one bit can affect the other bit (sometimes referred to complimentary bit disturb or CBD). Accordingly, it would be desirable to reduce feature sized so as to increase packing density while mitigating the adverse affects that may result there-from, such as Vt roll-off and CBD."}
{"text": "1. Field of the Invention\nThis technology relates to linearized power amplifiers (LPAs).\n2. Description of Related Art\nPrior LPAs have a variety of issues as follows.\n1. Prior LPA Efficiency and Linearity Issues:\nPrior linearized power amplifiers have several deficiencies including less than optimal efficiency, linearity, and performance over a voltage range and over a temperature range. One of the primary goals of any LPA is to provide a linear signal at as high a power level as possible. In a prior art LPA, the LPA current increases steadily with RF power out. At some high RF power level, clipping begins to occur to the RF signal thereby degrading linearity and efficiency at high RF power. One approach is to delay the onset of clipping, by reducing the bias current at high RF power levels to delay the onset of clipping.\n2. Prior LPA Performance Over Voltage Range Issue:\nPrior LPAs also tend to have significant variation in performance as power supply voltage level varies.\n3. Prior LPA Performance Over Temperature Range Issue:\nAnother shortcoming of prior LPAs is LPA performance over a temperature range. A prior art LPA with a fixed bias exhibits large changes in current/efficiency and linearity performance as the temperature changes from cold to room to hot.\n4. Other Prior LPA Issues:\nSome prior art LPAs attempt to improve these deficiencies by utilizing complex digitally modulated PA and digital pre-distortion schemes. However these types of designs are complex to design and implement, and tend to be large and consume much power. They are also usually silicon based technologies with inherently inferior RF characteristics, they cannot take advantage of the inherently superior RF characteristics of GaAs materials. Our novel LPA would be relatively simple and compact and able to be implemented on GaAs as a single die."}
{"text": "There are a very wide variety of requirements in sanitary fittings as to how the water guides to a shutoff valve and away from said shutoff valve are to be configured. Depending on the type and shape of the fitting, the water inlet may arrive at the shutoff valve from below or from the side, and also the outlet from the valve may adopt a very wide variety of directions."}
{"text": "1. Technical Field\nThe present disclosure relates to a discharge controller for controlling a discharge from a plurality of storage cells in a multiple cell battery to equalize remaining electricity storages in the storage cells.\n2. Related Art\nIn the vehicle such as Electric Vehicle (EV), Hybrid Electrical Vehicle (HEV), or the like, the multiple cell battery for feeding an electric power to a motor is installed, and a plurality of storage cells connected in series with each other are provided in the multiple cell battery.\nFIG. 23 is a block diagram showing a relationship between a multiple cell battery installed in the vehicle, a part of an electric drive system, and accessories. The vehicle shown in FIG. 23 is equipped with a multiple cell battery 10, an inverter 11, a motor 13, a DC-DC converter 15, a battery 17, and an electric motor 19 for a cooling fan (referred to as an “electric motor” hereinafter). An output voltage of the multiple cell battery 10 is a high voltage (e.g., 100 to 200 V), and an output voltage of the battery 17 is a voltage for the accessories (e.g., low voltage of 12V). The output voltage of the multiple cell battery 10 is converted from DC to AC by the inverter 11, and then is fed to the motor 13. Also, the output voltage of the multiple cell battery 10 is stepped down by the DC-DC converter 15, and then charges the battery 17. An electric power is supplied from the battery 17 to the electric motor 19. An airflow is generated by the cooling fan driven by the electric motor 19, and cools the multiple cell battery 10.\nThe multiple cell battery 10 shown in FIG. 23 includes a plurality of storage cells C1 to Cm (m is an integer of 2 or more) connected in series with each other (simply referred to as “cells” hereinafter), discharge switch units Cc1 to Ccm connected in parallel with the cells respectively, voltage detection units S1 to Sm connected in parallel with the cells respectively, and a control unit 21 for controlling respective discharge switch units. In this case, the discharge switch units Cc1 to Ccm, the voltage detection units S1 to Sm, and the control unit 21 are integrated and formed on an IC chip.\nEach discharge switch unit has a discharge resistor R and a switch Sw connected in series. The voltage detection units S1 to Sm detect a voltage across the cell (cell voltage) among the series-connected cells respectively. The control unit 21 controls ON/OFF of the switches in respective discharge switch units, in response to the result decided by the voltage detection units S1 to Sm.\nAlso, the multiple cell battery 10 has a charge control unit (not shown). The charge control unit performs the control to prevent an overcharge of each cell (overcharge preventing control), in charging the multiple cell battery 10. There is a variation in cell voltages of the cells depending on application conditions of the multiple cell battery 10, qualities of the cells, and the like. Therefore, the charge control unit executes the overcharge preventing control to meet the cell having the highest cell voltage. As the cells of the multiple cell battery 10, a lithium-ion battery, a nickel-hydrogen battery are employed, for example. In this case, this overcharge preventing control is particularly needed in employing the lithium-ion battery.\nIn this manner, the overcharge preventing control of the multiple cell battery 10 is applied based on the cell having the highest cell voltage. Therefore, as shown in FIG. 24, when a variation in cell voltages is increased due to repetition of the charge/discharge operation, a reduction in capacity of the multiple cell battery 10 is brought about. That is, the remaining electricity storage of the overall multiple cell battery 10 is restricted by the overcharge preventing control applied in charging the multiple cell battery 10. As a result, as indicated with a dotted line in FIG. 24, a capacity of the multiple cell battery 10 is lowered. When a sufficient electric power cannot be fed to the motor 13 because a capacity of the multiple cell battery 10 is lowered, either the cells of the multiple cell battery 10 must be equipped much more or the multiple cell battery 10 must be replaced with another multiple cell battery having a larger capacity.\nFor this reason, in the multiple cell battery 10 shown in FIG. 23, the control unit 21 controls individually the switches of respective discharge switch units such that cell voltages of the cells are kept at the same level. For example, as shown in FIG. 25, when the cell voltage of the cell C1 is higher than the cell voltages of the other cells C2 to Cm, the control unit 21 turns ON the switch Sw1 of the discharge switch unit Cc1 corresponding to the cell C1, and closes the circuit consisting of the cell C1 and a discharge resistor R1. At this time, a current flows to the discharge resistor R1 from the cell C1, and this current is converted into a heat by the discharge resistor R1. As a result, the cell voltage of the cell C1 is lowered, and the control unit 21 turns OFF a switch Sw1 when the cell voltage of the cell C1 is set at the same level as the cell voltages of the cells C2 to Cm. In this manner, a variation in cell voltages is reduced by equalizing the cell voltages of the cells, and therefore a reduction in capacity of the multiple cell battery 10 can be prevented (see e.g., JP-A-8-19188 and JP-A-2003-164069).\nIn the multiple cell battery 10 explained above, a heat is generated along with the equalization of the cell voltages of the cells. That is, in the multiple cell battery 10, an electric power accumulated in the cells is consumed wastefully by the equalization of the cell voltages of the cells. Also, as described above, the integrated circuit consisting of the discharge switch units Cc1 to Ccm, the voltage detection units S1 to Sm, and the control unit 21 is provided in the multiple cell battery 10. Since the constituent elements whose characteristics are varied depending on the ambient temperature are contained in the integrated circuit, it is advantageous that an amount of heat generated in the discharge resistors R1 to Rm should be made small.\nAccording to the Joule's law, an amount of heat generated in the resistor is proportional to “current2×resistance value”. Therefore, the resistor having a large resistance value is employed as the discharge resistors R1 to Rm. When a resistance value the discharge resistor is large, a discharge current becomes small, so that a quantity of heat generated by the discharge resistor is reduced. However, when the discharge current is small, a time required for the equalization of the cell voltages of the cells is prolonged. As a result, the resistance value of the discharge resistors R1 to Rm is decided in consideration of the balance between a quantity of heat generated by the discharge resistor and an equalizing time for the cell voltages. In any event, a heat is generated from the discharge resistors R1 to Rm along with the equalization of the cell voltages of the cells. Therefore, high-level heat measures such as thermal diffusion, heat resistance, etc. must be applied to the multiple cell battery 10 such that the integrated circuit is not affected by such heat."}
{"text": "1. Field of the Invention\nThe invention relates to a curable composition and a cured polymer product.\n2. Description of the Related Art\nThe ion exchange membrane is used in electrodeionization (EDI), continuous electrodeionization (CEDI), electrodialysis (ED), and electrodialysis reversal (EDR). The ion exchange membrane is used not only for general usage but also for medical use, and recently the ion exchange membrane has also been used in a solid polymer electrolyte-type fuel cell.\nThe electrodeionization (EDI) is a water treatment process for removing ions from aqueous liquids by using an ion exchange membrane and a potential in order to achieve ion transport. EDI differs from other water purification technologies, such as conventional ion exchange, in that it is does not require the use of chemicals such as acids or caustic soda, and can be used to produce ultra pure water. The electrodialysis (ED) and the electrodialysis reversal (EDR) are electrochemical separation processes for removing ions and the like from water and other fluids.\nRecently, the ion exchange membrane has been generally manufactured in the method of thermally polymerizing a styrene-based monomer and introducing an ion exchangeable substituent (quaternary ammonium salt group or sulfo group) by substitution reaction or the like (for example, see JP2006-225459A). Recently, research for improving a film provided with ionic groups and crosslinking bonds at high density has been actively conducted by using a monomer (also referred to as a charged crosslinker or an ionic crosslinking agent) also having an ionic group and a crosslinking group (2 or more polymerizable groups). For example, mainly in a polymer, an anion exchange membrane (for example, see WO2013/011273A) having a cationic group such as quaternary ammonium or a cation exchange membrane (for example, see WO2013/011272A, JP2007-305371A, and JP2011-189223A) having an anionic group such as a sulfo group is suggested."}
{"text": "1. Field of the Invention\nThis invention relates to an improvement of a data reproducing device which was previously filed by the applicant of this invention in U.S. Patent Office on Dec. 14, 1993, and is disclosed in U.S. patent application Ser. No. 08/165,885, and more particularly to a data reproducing device which is suitably used to reproduce both of video data and audio data recorded on an optical disc in a time division multiplexing (TDM) imode while separating and decoding these data.\n2. Description of Related Art\nFIG. 1 is a block diagram showing the data reproducing device which is disclosed in U.S. patent application Ser. No. 08/165,885 (FIG. 6).\nThe data reproducing device has a disc driver 1 for reproducing data from an optical disc, on which video data and audio data are recorded in a time division multiplexing (TDM) mode, a decoding unit 2 to which reproduction data output from the driver 1 is supplied to decode the input data, and an ECC circuit 3 for receiving data output from the decoding unit 2 to detect and correct errors of the data. The output of the ECC circuit 3 is supplied to each of a ring buffer 4 and an address extracting circuit 31. The ring buffer 4 serves to accumulate a predetermined amount of data input thereto, and then output the accumulated data to a multiplexed data separating unit 5.\nThe multiplexed data separating unit 5 has a data separation circuit 21 for separating each of video data, audio data, an SCR (system clock reference), a DTS (decoding time stamp) for video data (hereinafter referred to as \"DTSV\"), and a DTS for audio data (hereinafter referred to as \"DTSA\") from data supplied from the ring buffer 4.\nA data format for the data to be supplied to multiplexed data separation unit 5 is defined as shown in FIG. 2, for example. The format is defined as a multiplexed bit stream of MPEG (ISO11172). As shown in FIG. 2, the multiplexed bit stream comprises one or more packs (PACK), and each pack comprises one or more packet (PACKET). At the head of the pack is disposed a pack header (PACK HEADER), and at the packheader are disposed a pack start code (PACK START CODE) representing a start point of the pack, SCR and MUX RATE. The SCR represents a time when its final byte is input to the multiplexed data separation unit 5 (a time when a demultiplexing operation is started). The MUX RATE represents a transmission rate.\nIn an embodiment shown in FIG. 2, a video packet (VIDEO PACKET) and an audio packet (AUDIO PACKET) are successively disposed subsequently to the pack header. A packet header (PACKET HEADER) is disposed at the head of each of these video and audio packets, and the packet header has a video packet start code (VIDEO PACKET START CODE) or audio packet start code (AUDIO PACKET START CODE) representing a start point of the video or audio packet, and a DTSV or DTSA representing a start time of a decoding operation of the video or audio data. The video data (VIDEO DATA) or audio data (AUDIO DATA) are disposed subsequently to each packet header.\nThe timing data (time information) such as SCR, DTS (DTSV or DTSA) are represented with a count value of 90 kHz frequency clocks, and have 33-bit significant digit.\nThe video data which are separated in the data separation circuit 21 are supplied to a video code buffer 6 (FIFO), and the audio data are supplied to an audio code buffer 8 (FIFO). The SCR is supplied to and stored in an STC register 26. The STC register 26 serves to count clocks of 90 kHz frequency output from a clock generating circuit 27 and increment its storage value to generate an STC (system time clock).\nThe DTSV and DTSA which are separated by the data separation circuit 21 are supplied to and stored in a DTSV register 22 and a DTSA register 24 respectively. The data stored in the DTSV register 22 and the DTSA register 24 are supplied to comparators 23 and 25 respectively, and these data are compared with the STC output from the STC register 26.\nThe data reproducing device shown in FIG. 1 also has a control circuit 28 which comprises a CPU or the like and controls a reproducing operation on the basis of an instruction which is input in correspondence with an user's manipulation of an input unit 29, and an address storing circuit 30 for storing a reproduction start address and a reproduction end address which are input through the input unit 29.\nThe video data stored in the video code buffer 6 are read out and supplied to a video decoder 7, and decoded by the video decoder 7 to generate video signals. The generated video signals are output to a circuit (not shown). The video decoder 7 is supplied With a video decode start signal output from the comparator 23.\nLikewise, the data output from the audio code buffer 8 are supplied to an audio decoder 9, and decoded therein. The audio decoder 9 is also supplied with an audio decode start signal output from the comparator 25.\nNext, the operation of the data reproducing device as described above will be described with reference to FIG. 3. First, through manipulation of the input unit 29, a reproduction start address, a reproduction end address and a reproduction start are instructed to the control circuit 28. The reproduction start address and the reproduction end address are stored in the address storing circuit 30. At this time, the control circuit 28 outputs an instruction to the driver 1 to reproduce data recorded on an optical disc which is mounted in the driver 1. The reproduction data output from the driver 1 are supplied to the decoding unit 2 and decoded therein, and then supplied to the ECC circuit 3 in which the data are subjected to an error detection and correction processing.\nThe data which have been subjected to the error detection and correction processing in the ECC circuit 3 are supplied to the address extraction circuit 31 to read out an address disposed at a predetermined position in the data. The read-out address is supplied to the control circuit 28. The control circuit 28 is on standby until the read-out address supplied thereto is coincident with the reproduction start address from the address storing circuit 30, and if both of the addresses are coincident with each other, the control circuit 28 generates and outputs an accumulation start command to the ring buffer 4. Through this operation, the reproduction data output from the driver 1 are accumulated in the ring buffer 4. Thereafter, the data stored in the ring buffer 4 are read out from the ring buffer 4, and supplied to the data separation circuit 21 of the multiplexed data separation unit 5.\nThe data separation circuit 21 is controlled by the control circuit 28 to separate the video data and the audio data from the data supplied from the ring buffer 4. These video data and audio data are supplied to the video code buffer 6 and the audio code buffer 8, respectively. In addition, the data separation circuit 21 separates the SCR, the DTSV and the DTSA from the data, and supplied to and stored in the STC register 26, the DTSV register 22 and the DTSA register 24, respectively.\nThe STC register 26 stores the SCR, and subsequently it counts the clocks output from the clock generating circuit 27 to increment its storage value (SCR) in response to each clock. The storage value of the STC register 26 is supplied to the comparators 23 and 25 as an inner time (STC).\nThe DTSV register 22 keeps a DTSV which is first supplied thereto after the reproduction is started by the driver 1. This DTSV corresponds to a decode start time of the data of a picture at the head of the video data which are separated by the data separation circuit 21 and supplied to and stored in the video code buffer 6.\nLikewise, the DTSA register 24 keeps a DTSA which is first supplied thereto after the reproduction is started. This DTSA corresponds to a decode start time of the data of a decode unit at the head of the audio data which are stored in the audio code buffer 8.\nThe SCR corresponds to a time when the data are supplied from the ring buffer 4 to the multiplexed data separation unit 5 and the demultiplexing operation is started (a time at the input time point of the input data). That is, it corresponds to a time t1 in FIG. 3. The STC register 26 serves to input a time data (current time) from the time t1 to one inputs of the comparators 23 and 25.\nThe DTSV register 22 supplies the other input of the comparator 23 with the time DTSV when the video decoder 7 starts its decoding operation. The comparator 23 outputs the video decode start signal to the video decoder 7 when the current time (STC) output from the STC register 26 is coincident with the decode start time output from the DTSV register 22 (time t2 in FIG. 3). The video decoder 7 reads out one frame of video data which have been written in the video code buffer 6 when supplied with the video decode start signal, and starts its decoding operation.\nIn FIG. 3, a straight line A represents a writing state of data into the video code buffer 6 (its slope represents a write-in transfer rate), and a polygonal line B represents a read-out state of data from the video code buffer 6 by the video decoder 7. Accordingly, data remains in an area represented by a shadow in the video code buffer 6. The storage capacity of the video code buffer 6 is represented by the distance in a vertical direction between lines A and C.\nUpon input of the video code start signal, the video decoder 7 starts its decode operation to generate a video vertical synchronizing signal at the time when the decoding operation is completed, that is, at the time when a video decode delay time (VIDEO DECODE DELAY) elapses after the start of the decode operation, and output video signals subsequently to this generation. Accordingly, a display is started when the video decode delay time elapses after the start of the decoding operation.\nLikewise, the comparator 25 outputs the audio decode start signal when the current time (STC) output from the STC register 26 is coincident with the decode start time (DTSA) of the audio data output from the DTSA register 24. When the audio decode start signal is input, the audio decoder 9 reads out from the audio code buffer 8 data whose amount corresponds to a decode unit and starts its decode processing to generate audio signals, and outputs the generated audio signals to a circuit (not shown).\nDuring the reproducing operation of the driver 1, addresses being read out are supplied from the address extraction circuit 31 to the control circuit 28 at all times. The control circuit 28 compares the address of data being read out and supplied thereto with the reproduction end address stored in the address storing circuit 30, and if both addresses are coincident with each other, it ends the read-out operation from the driver 1.\nSince the data remain in the ring buffer 4 and the code buffers 6 and 8, the actual end time of the reproducing operation is the time when there exists no data to be input to the video decoder 7 and the audio decoder 9. In this case, the control circuit 28 finishes the reproducing operation when detecting the absence of the data to be input to the video decoder 7 and the audio decoder 9.\nNow, it is considered to continuously reproduce any discontinuous two points on an optical disc having video signals recorded thereon while compressed in the MPEG system. In other words, an edition of the reproducing operation is now considered.\nFIG. 4 is a timing chart for a continuous reproducing operation of a predetermined range of data on the optical disc. Here, sectioning points for an edition are represented by G, H and I, and pictures before and after the sectioning points G, H and I are represented by (a,b), (c,d) and (e, f), respectively.\nIn FIG. 4, a straight line D represents a write-in state of data into the video code buffer 6, and its slope represents a write-in transfer rate. A polygonal line E represents a data read-out state from the video code buffer 6 by the video decoder 7. Accordingly, data remain in an area represented by a shadow in the video code buffer 6. The storage capacity of the video code buffer 6 is represented by the distance in a vertical direction between the lines D and F.\nIt is now assumed that the reproduction point is jumped from the point G to the point H in FIG. 4. If a continuous reproduction is performed with no jumping operation, a picture b is reproduced subsequently to a picture a. However, if the jumping operation is performed, a picture d is reproduced subsequently to the picture a. Assuming a time required for the jump from the point G to the point H to be zero, the write-in and read-out operations of the data of the video code buffer 6 when the jumping operation as described above is performed are carried out as shown in FIG. 5.\nThat is, when no jumping operation is performed, at the point G, the write-in operation of the picture b and subsequent pictures are performed for the video code buffer 6 at a time t11, and this write-in operation is continued until a time t12. At the time t12, the data of the picture b is read out from the video code buffer 6, and supplied to the video decoder 7 to decode the data. The read-out of the data of the picture b at the time t12 is started after a time corresponding to one frame elapses from the read-out of the data of the picture a just before the picture b. The time period Tb between the time t11 and the time t12 is a start up delay of the picture b, and it corresponds to a time from the input of the data of the picture into the video code buffer 6 till the decoding of the data. In the MPEG system, an image compression processing is performed in accordance with complication of an image of each picture, so that a data amount becomes larger as an image is complicated while the data amount becomes smaller as an image is simple. Therefore, a different start up delay is basically adopted every picture in the MPEG system.\nIn the above case, the start up delay Tb of the picture b to be reproduced subsequently to the picture a when no jump is performed at the point G becomes smaller than the start up delay Td of the picture d to be reproduced subsequently to the picture a through the jumping operation from the point G to the point H.\nAccordingly, as shown in FIG. 5, just before the jump is performed at the point G, the video decoder 7 is about to read out and decode the data of the next picture b from the video code buffer 6 at the timing when the data of the picture a written in the video code buffer 6 is decoded and a time corresponding to one frame elapses from the above decoding. However, in this case, since the jump is performed, the write-in operation of the data of the picture d into the video code buffer 6 is started at the time t11. Therefore, when the video decoder 7 is about to read out and decode the data of a next picture at the time t12 when the start up delay Tb of the picture b elapses from the time t11, one frame data of the picture d has not yet been written in the video code buffer 6 because the start up delay Td of the picture d is larger than the start up delay Tb of the picture b. Consequently, if the video decoder 7 reads out one frame data from the video code buffer 6, the video code buffer 6 falls into an underflow state.\nOn the other hand, in a case where the jump is performed from the point H to the point I in FIG. 4, the start up delay Td of the picture d to be reproduced subsequently to a picture c when no jump is performed at the point H becomes larger than the start up delay Tf of a picture f to be reproduced subsequently to the picture c through the jump operation from the point H to the point I. Consequently, the operation of the video decoder 7 when the jump is performed from the point H to the point I is performed as shown in FIG. 6 on the assumption that the time required for the jump operation is equal to zero as described above.\nIn this case, through the jump from the point H to the point I at a time t13, the write-in operation of the data of the picture f into the video code buffer 6 is started at the time t13. However, the start up delay Tf of the picture f is smaller than the start up delay Td of the picture d which would be reproduced if no jump is performed at the point H, so that the data of the picture f will have been stored in the video code buffer 6 for a time longer than the original start up delay Tf of the picture f.\nThat is, the video decoder 7 cannot start the decoding of the data of the picture f from the time when it reads out the data of the picture c until a time t32 when a time corresponding to one frame elapses from the above read-out time (this time corresponds to the end time of the start up delay Td of the picture d). As a result, the decode timing for data which contain the picture f data and subsequent picture data (read-out timing from the video code buffer 6) is delayed by the difference between the start up delay Td of the picture d and the start up delay Tf of the picture f (Td-Tf). Therefore, as shown in FIG. 6, the picture f data and the subsequent data are accumulated, and the video code buffer may fall into an overflow state.\nIn order to prevent the underflow as described above, the decode timing of the video decoder 7 must be delayed from the time t12 in FIG. 5. Likewise, in order to prevent the overflow, the write-in of the data into the video code buffer 6 from the point I when the jump is performed at the point H is required to be delayed by the difference (Td-Tf) of the start up delays of the pictures d and f.\nHowever, the device shown in FIG. 1 has no element for adjusting the start up delay before and after an edition point. Therefore, when an editing operation is instructed, the overall reproducing operation is temporarily stopped at the reproduction point at which the jump is instructed. Thereafter, the reproducing operation is re-started at a jumped reproduction point.\nAs a result, for example when the jump is performed from the point G to the point H, the overall reproducing operation of the driver 1 is temporarily stopped at the point G (time t11) as shown in FIG. 8. However, since the data of the picture a has been already written in the video code buffer 6, the video decoder 7 reads out the data of the picture a at a time t41 to decode and output the data. After the decode processing of the data of the picture a is completed, the control circuit 28 controls the driver 1 again to restart the reproducing operation from the point H (time t42). Through this operation, at a time t43 when the start up delay Td of the picture d elapses from the time t42, the data are read out from the video code buffer 6, and then decoded in the video decoder 7.\nIn a case where the jump is performed from the point H to the point I, the overall reproducing operation of the driver 1 is temporarily stopped at the time t13 (when arriving at the point H) as shown in FIG. 9. Thereafter, the data which have been already written in the video code buffer 6 are decoded in the video decoder 7. The decode processing is completed at a time t51.\nAs described above, when the decoding operation of the data which have been already stored in the video code buffer 6 is completed, the driver 1 is controlled again, and the reproduction is started from the point I at the time t52 to write the data in the video code buffer 6. Thereafter, at the time t53 when the start up delay Tf of the picture f elapses from the time t52, the decode of the picture f is started.\nIn the device shown in FIG. 1, when the reproduction point is discontinuously shifted as described above, the reproducing operation is temporarily finished at a predetermined reproduction point, and then the reproducing operation is restarted from a shifted reproduction point. Therefore, there exists a problem of a data lack time of reproduction data, that is, a time in which the reproduction data are lacking (the period from the time t41 till the time t43 in FIG. 8 or the period from the time t51 till the time t53 in FIG. 9)."}
{"text": "1. Field of the Invention\nThe field of invention relates to rescue apparatus, and more particularly pertains to a new and improved rescue beacon apparatus wherein the same is arranged for selective inflation and disposition above an individual requiring assistance.\n2. Description of the Prior Art\nThe prior art has utilized inflatable gas balloons wherein U.S. Pat. No. 4,794,498 to Neumeier sets forth a gas balloon utilizing an illumination member therewithin for use in an advertising situation, wherein a housing mounting a battery container is secured relative to a lower end portion of the balloon.\nU.S. Pat. No. 4,901,664 to Labrecque sets forth a survival kit utilizing an inflatable balloon arranged for projection relative to an individual supporting a holder mounting a reel therewithin.\nU.S. Pat. No. 4,787,575 to Stewart sets forth a signal balloon mounted relative to a base container that is ground supported and includes a quantity of tether line for securing the balloon relative to the housing.\nU.S. Pat. No. 4,944,242 to Russell sets forth a rescue balloon including an inflation cylinder and a tether line for securing the balloon relative to a ground supported container.\nAs such, it may be appreciated that there continues to be a need for a new and improved rescue beacon apparatus as set forth by the instant invention which addresses both the problems of ease of use as well as effectiveness in construction and in this respect, the present invention substantially fulfills this need."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a plug seedling extractor to facilitate the removal or extraction of seedlings from plug seedling trays. After germination, transplantation of the tiny fragile seedlings from the seedling tray to a larger growing environment is necessary. The present invention quickly and efficiently extracts the seedlings from the plug seedling tray with minimal damage to the plug of soil, the seedling, and its root system.\n2. Description of Related Art\nTraube in British Patent No. 341,819 dated Jan. 22, 1931 discloses a device for molding seedling plugs and extracting the molded seedling plugs from the mold. The device is impractical because the seedlings are grown apart from the device, molded in the device, and then presumably transplanted into another device immediately after molding. The device is too bulky for mass marketing or shipping.\nBlackmore, et al, in U.S. Pat. No. 3,799,078 generally discloses a method of transplanting plug seedlings from a plug seedling tray to larger containers. The method involves extractor pins with concave openings at the tips of the pins. The extractor pins push the plug seedlings downwardly through the plug seedling tray. There is a great likelihood of damage to the foliage of the seedling if it is not extracted from the plug seedling tray relatively soon after germination.\nBlackmore, Jr., in U.S. Pat. No. 4,197,674 discloses a plug seedling tray for use with extractor pins that pass upwardly through the bottom of the seedling tray to extract the plug seedlings, but no particulars of the device used with the plug seedling tray are provided.\nAccordingly, no device is disclosed in prior art suitable for use with a plurality of variously-sized plug seedling trays for extracting the seedlings with minimal damage to the seedling itself and to the plug containing the root system."}
{"text": "1. Field of the Invention\nThis invention relates to an image formation apparatus such as a copying apparatus or a recording apparatus having a mechanism for manually or automatically feeding sheets.\n2. Description of the Prior Art\nThere are known apparatus in which a sheet is fed from a place such as a cassette containing a number of sheets therein and exposure or image transfer is applied to the sheet to obtain a copy image or a sheet is manually fed to obtain a copy image thereon.\nIt is particularly convenient to manually feed a sheet along another path and preferentially obtain a copy image on the sheet which is different in size or paper quality from the sheets contained in the cassette.\nHowever, provision of independent and exclusive paths for manual supply and automatic supply from a cassette which lead to the image formation station would result in a bulky and complicated construction of the copying apparatus.\nAlso, feeding sheets by the use of a common path would cause jam of the sheets or apparatus trouble which would be inconvenient to the operator.\nIn a case where the cassette process sequence is controlled by the same control circuit in accordance with various signals from the cassette, the process sequence concerned with the manual supply may also be controlled by those signals and this may cause a process trouble or cause formation of an incomplete image on the manually supplied sheet.\nIt would also be convenient to make it possible to form an image or manually supplied sheets of different sizes. However, when a sheet of a width smaller than the width of the insertion bed is inserted, the sheet may be obliquely inserted with a result that an image is formed at an unsuitable position on the sheet or the sheet jams.\nIn a case where a sheet is supported and fed so as to be inserted into the apparatus and then a copy start button is operated, the operator is required to use both of his hands and this means very poor operability. Moreover, it is very difficult to insert a sheet straight feed and it is often the case that the sheet is fed obliquely.\nParticularly, in a transfer type copying apparatus, where the sequence control is effected by detecting a manually supplied sheet, it is necessary to select the position of the sheet detector to an optimal position so as not to harm the image transfer registration. This imposes a limitation on the compactness of the copying apparatus.\nIn a transfer type copying apparatus, if the path leading to the image transfer station is long, sheet jam or image transfer misregistration is liable to occur. For example, there are apparatus in which a sheet is stopped by stopped register rollers and image transfer registration is effected by driving a roller in synchronism with the image on the drum. In this case, the sheet is caused to strike against the register rollers to provide a flexure of the sheet, but if the path to the register rollers is long, an appropriate flexure may not be provided sometimes. This may lead to poor registration accuracy or occurrence of sheet jam.\nAlso, in a transfer type copying apparatus, if the photosensitive drum is made small, the time required for the leading end of the image on the drum to reach the image transfer station from the start of copying becomes so much shorter that there is not enough time allowed for good image transfer registration to be obtained, and this may result in poor registration accuracy.\nFurther, in a transfer type copying apparatus, an image has been formed on the photosensitive drum by reciprocal movement of an original platen to make the apparatus compact and to obtain copies of a relatively large size. However, a complicated mechanism is necessary to move the platen and automatically return it and thus, low cost and compactness of the apparatus have been difficult to provide.\nAlso, the compactness of the copying apparatus has been limited due to the size of its photosensitive drum and the length of the optical path of the original image relative to the drum."}
{"text": "1. Field of the Invention\nThis invention relates to cartons and more particularly to the combination of a shipping carton and merchandise fixture for attachment to a vertically oriented structural member, such as a conventional slotted upright or standard for shelving, a wire rack, a slat wall, a peg board, etc.\n2. Description of Related Art\nVarious structures for displaying merchandise are commonly found in commercial enterprises and several are the subject of patents. Many of such structures make use of a horizontally disposed member arranged to be mounted on a vertical surface to support, e.g., suspend, merchandise for sale or other purposes. See for example, U.S. Pat. No. 6,202,866 (Shea), U.S. Pat. No. 6,199,706 (Shea) and U.S. Pat. No. 6,070,747 (Shea).\nU.S. Pat. No. 4,779,720 (Mandelbaum) discloses a combination shipping and display device for neckties and the like. The device basically comprises a generally rectangular, open faced inner carton. That carton includes an assembly of hanger rods supported from an assembly of horizontal and vertical rods. The assembly of hanger rods and the associated horizontal and vertical rods are located and secured within the inner carton. The top of the inner carton folds backward and includes a flap affording an advertising panel and a tongue retained between the carton back panel and a top horizontal rod member in the display mode. A retainer assembly is also provided and includes spaced U-shaped apertures and baffles for laterally supporting columns of ties on the hangers. The inner carton and retainer assembly are inserted into a conventional outer carton for shipping. The assembly of hanger rods and the inner carton in which they are mounted, are arranged to be connected via a hook to a mesh panel, peg board or other vertical structure to support the inner carton holding the ties on that vertical structure.\nU.S. Pat. No. 4,842,131 (Mandelbaum) discloses a combination shipping and display device for neckties and the like that is similar to the device of U.S. Pat. No. 4,779,720 (Mandelbaum), but which includes a hanger support structure. In particular, the hanger support structure is similar to the assembly of the hanger rods and associated horizontal and vertical rods, but is in the form of an outward facing U including a pair of end strips with holes for receipt of “Christmas tree” fasteners to secure the hanger support structure in the inner carton. The hangar support structure and the inner carton in which it is mounted, are arranged to be connected via a hook to a mesh panel, peg board or other vertical structure to support the inner carton holding the ties on that vertical structure.\nU.S. Pat. No. 5,249,668 (Fenton et al.) discloses a revolvable rack for neckties that can be either temporarily or permanently hooked or fastened to a rod or a store fixture, and a container for shipping the rack while it is fully loaded with ties. The rack is easily rotated by hand, has a relatively small turning radius and has the capacity to hold several dozen ties. Once the tie rack is placed in the container, the container serves to both store and protect the ties in an organized and out-of-the-way manner. The tie rack and container combination can be hung together in an existing space. The container allows the fully loaded tie rack to be quickly and easily lifted out of the container in one smooth motion, and directly hung on a rod or on a merchandising fixture in a store without the use of special tools or equipment, such that the ties are readily and pleasingly displayed.\nOther cartons including or adapted to make use hangers for garments or other merchandise are disclosed in the following U.S. Pat. No. 2,796,977 (Divine), U.S. Pat. No. 3,659,704 (Collura et al.), U.S. Pat. No. 3,987,898 (Crane), U.S. Pat. No. 4,098,399 (Behtune et al.), U.S. Pat. No. 4,576,280 (Dove et al.), U.S. Pat. No. 4,693,369 (Lagin) and U.S. Pat. No. 6,155,415 (Runyan).\nWhile the devices of the aforementioned prior patents appear generally suitable for their intended purposes, they still leave something to be desired from one or more standpoints, such as simplicity of construction, ease of opening the container, ease of use, ease of disassembly of the container, ease of assembly of the container's fixture components, ease of carrying the container to effect the mounting of the fixture onto a vertical support surface and the ability of the fixture to be mounted on various types of vertical support surfaces to project out from those surfaces in either of two opposite directions.\nIn our copending U.S. patent application Ser. No. 10/305,501, filed on Nov. 27, 2002, entitled Ambidextrous Merchandise Fixture And Method of Displaying Merchandise Therefrom, which is assigned to the same assignee as this invention and whose disclosure is incorporated by reference herein, there is disclosed a fixture and method of use for supporting merchandise, e.g., plural garment belts, on a vertical support structure. That fixture is arranged to be mounted any one of various types of conventional vertical oriented support structures, such as a conventional slotted upright, a wire rack, a slat wall and a peg board, with the fixture projecting out from any of those structures in either of two opposite directions. To that end the fixture includes a bracket and an elongated display member having a pair of ends. The bracket includes a section for snap connection to either end of the display member and a connector that is configured to engage the vertically oriented support structure. The display member includes merchandise holders, e.g., plural prong hangers, for holding the merchandise thereon, e.g., suspending the merchandise therefrom.\nThe subject invention makes use of the fixture of our aforementioned patent application in an integrated system, i.e., a system including a shipping carton and the merchandise, that is simple in construction, low in cost and easy to use."}
{"text": "The present invention relates to compositions, in either liquid or granular form, for use in laundry applications, wherein the compositions comprise certain amino acid based polymer, oligomer or copolymer materials which impart appearance and integrity benefits to fabrics and textiles laundered in washing solutions formed from such compositions.\nIt is, of course, well known that alternating cycles of using and laundering fabrics and textiles, such as articles of worn clothing and apparel, will inevitably adversely affect the appearance and integrity of the fabric and textile items so used and laundered. Fabrics and textiles simply wear out over time and with use. Laundering of fabrics and textiles is necessary to remove soils and stains which accumulate therein and thereon during ordinary use. However, the laundering operation itself, over many cycles, can accentuate and contribute to the deterioration of the integrity and the appearance of such fabrics and textiles.\nDeterioration of fabric integrity and appearance can manifest itself in several ways. Short fibers are dislodged from woven and knit fabric/textile structures by the mechanical action of laundering. These dislodged fibers may form lint, fuzz or xe2x80x9cpillsxe2x80x9d which are visible on the surface of fabrics and diminish the appearance of newness of the fabric. Further, repeated laundering of fabrics and textiles, especially with bleach-containing laundry products, can remove dye from fabrics and textiles and impart a faded, worn out appearance as a result of diminished color intensity, and in many cases, as a result of changes in hues or shades of color.\nGiven the foregoing, there is clearly an ongoing need to identify materials which could be added to laundry detergent products that would associate themselves with the fibers of the fabrics and textiles laundered using such detergent products and thereby reduce or minimize the tendency of the laundered fabric/textiles to deteriorate in appearance. Any such detergent product additive material should, of course, be able to benefit fabric appearance and integrity without unduly interfering with the ability of the laundry detergent to perform its fabric cleaning function. The present invention is directed to the use of amino acid based polymer, oligomer or copolymer materials in laundry applications which perform in this desired manner.\nAmino acid based polymer, oligomer or copolymer materials which are suitable for use in laundry operations and provide the desired fabric appearance and integrity benefits can be characterized by the following general formula: \nwherein the polymer, oligomer, or copolymer contains at least about 5 mole %, preferably at least about 10 mole %, more preferably from about 20 mole %, and most preferably at least about 40 mole %, of a basic amino acid;\neach R1 is selected from the group consisting of H, C1-C18 saturated or unsaturated, branched or linear alkyl, C2-C18 saturated or unsaturated, branched or linear hydroxyalkyl, C3-C8 cycloalkyl, C6-C18 aryl, and C7-C18 alkylaryl;\neach R2 is independently selected from the group consisting of H, NH2, \neach R3 is independently selected from the group consisting of OH, OM, N(R1)2, \neach R5 is independently selected from the group consisting of C1-C12 linear or branched alkylene, cyclic alkylene, C2-C12 linear oxa-substituted alkylene, C2-C12 branched oxa-substituted alkylene, C3-C12 cyclic oxa-substituted alkylene, and \nxe2x80x83wherein:\neach x is independently from 0 to about 200;\neach y is independently from 0 to about 10, preferably y is 0, 3, or 4, and most\npreferably y is 4;\neach z is independently from 1 to about 7;\neach a is independently from about 1 to about 12;\nM is selected from compatible cations; and\nprovided that:\nthe sum of all x\"\"s is from 2 to about 200, preferably from about 3 to about 150, more preferably from about 5 to about 120, and most preferably from about 5 to about 100;\nany basic amine site on the polymer, oligomer, or copolymer may be optionally protonated, alkylated, or quaternized with groups selected from the group consisting of H, CH3, alkyl, hydroxyalkyl, benzyl and mixtures thereof;\nany amine site may be optionally alkoxylated; and\nwhen two R1 groups are attached to a common nitrogen the two R1s may form a cyclic structure selected from the group consisting of C5-C8 alkylene, and C4-C7 alkyleneoxyalkylene.\nThe amino acid based polymer, oligomer or copolymer materials defined above can be used as a washing solution additive in either granular or liquid form. Alternatively, they can be admixed to granular detergents, dissolved in liquid detergent compositions or added to a fabric softening composition. Preferably the fabric treatment compositions of this invention comprise from about 0.1% to about 10%, preferably from about 0.2% to about 8%, more preferably from about 0.3% to about 6%, and most preferably from about 0.4% to about 5%, by weight of a mixture of the amino acid based polymers, oligomers or copolymers defiened by the general formula above. The forgoing description of uses for the amino acid based fabric treatment materials defined herein are intended to be exemplary and other uses will be apparent to those skilled in the art and are intended to be within the scope of the present invention.\nThe laundry detergent compositions herein comprise from about 1% to 80% by weight of a detersive surfactant, from about 0.1% to 80% by weight of an organic or inorganic detergency builder and from about 0.1% to 5% by weight of the amino acid based fabric treatment materials of the present invention. The detersive surfactant and detergency builder materials can be any of those useful in conventional laundry detergent products.\nAqueous solutions of the amino acid based polymer, oligomer or copolymer materials of the subject invention comprise from about 0.1% to 50% by weight of the amino acid based fabric treatment materials dissolved in water and other ingredients such as stabilizers and pH adjusters.\nIn its method aspect, the present invention relates to the laundering or treating of fabrics and textiles in aqueous washing or treating solutions formed from effective amounts of the detergent compositions described herein, or formed from the individual components of such compositions. Laundering of fabrics and textiles in such washing solutions, followed by rinsing and drying, imparts fabric appearance benefits to the fabric and textile articles so treated. Such benefits can include improved overall appearance, pill/fuzz reduction, antifading, improved abrasion resistance, and/or enhanced softness.\nAs noted, when fabric or textiles are laundered in wash solutions which comprise the amino acid based polymer, oligomer or copolymer materials of the present invention fabric appearance and integrity are enhanced. The amino acid based fabric treatment materials can be added to wash solutions by incorporating them into a detergent composition, a fabric softener or by adding them separately to the washing solution. The amino acid based fabric treatment materials are described herein primarily as liquid or granular detergent additives but the present invention is not meant to be so limited. The amino acid based fabric treatment materials, detergent composition components, optional ingredients for such compositions and methods of using such compositions, are described in detail below. All percentages are by weight unless other specified.\nA) Amino Acid Based Polymer Oligomer or Copolymer Materials\nThe essential component of the compositions of the present invention comprises one or more amino acid based polymer, oligomer or copolymer. Such materials have been found to impart a number of appearance benefits to fabrics and textiles laundered in aqueous washing solutions formed from detergent compositions which contain such amino acid based fabric treatment materials. Such fabric appearance benefits can include, for example, improved overall appearance of the laundered fabrics, reduction of the formation of pills and fuzz, protection against color fading, improved abrasion resistance, etc. The amino acid based fabric treatment materials used in the compositions and methods herein can provide such fabric appearance benefits with acceptably little or no loss in cleaning performance provided by the laundry detergent compositions into which such materials are incorporated.\nThe amino acid based polymer, oligomer or copolymer component of the compositions herein may comprise combinations of these amino acid based materials. For example, a mixture of lysine and hexamethylenediamine condensates can be combined with a mixture of lysine and octylamine condensates to achieve the desired fabric treatment results. Moreover, the molecular weight of amino acid based fabric treatment materials can vary within the mixture as is illustrated in Example I below.\nAs will be apparent to those skilled in the art, an oligomer is a molecule consisting of only a few monomer units while polymers comprise considerably more monomer units. For the present invention, oligomers are defined as molecules having an average molecular weight below about 1,000 and polymers are molecules having an average molecular weight of greater than about 1,000. Copolymers are polymers or oligomers wherein two or more dissimilar monomers have been simultaneously or sequentially polymerized. Copolymers of the present invention can include, for example, polymers or oligomers polymerized from a mixture of a primary amino acid based monomer, e.g., lysine, and a secondary amino acid monomer, e.g., tryptophan.\nThe amino acid based fabric treatment component of the detergent compositions herein will generally comprise from about 0.1% to about 10%, preferably from about 0.2% to about 8%, more preferably from about 0.3% to about 6%, and most preferably from about 0.4% to about 5%, by weight of a mixture of the amino acid based polymers, oligomers or copolymers defiened by the general formula below. But when used as a washing solution additive, i.e. when the amino acid based fabric treatment component is not incorporated into a detergent composition, the concentration of the amino acid based component can comprise from about 0.1% to about 50% by weight of the additive material.\nOne suitable group of amino acid based polymer, oligomer or copolymer materials for use herein is characterized by the following formula: \nwherein the polymer, oligomer, or copolymer contains at least about 5 mole %, preferably at least about 10 mole %, more preferably from about 20 mole %, and most preferably at least about 40 mole %, of a basic amino acid;\neach R1 is selected from the group consisting of H, C1-C18 saturated or unsaturated, branched or linear alkyl, C2-C18 saturated or unsaturated, branched or linear hydroxyalkyl, C3-C8 cycloalkyl, C6-C18 aryl, and C7-C18 alkylaryl;\neach R2 is independently selected from the group consisting of H, NH2, \neach R3 is independently selected from the group consisting of OH, OM, N(R1)2, \neach R5 is independently selected from the group consisting of C1-C12 linear or branched alkylene, cyclic alkylene, C2-C12 linear oxa-substituted alkylene, C2-C12 branched oxa-substituted alkylene, C3-C12 cyclic oxa-substituted alkylene, and \nxe2x80x83wherein:\neach x is independently from 0 to about 200;\neach y is independently from 0 to about 10;\neach z is independently from 1 to about 7;\neach a is independently from about 1 to about 12;\nM is selected from compatible cations; and\nprovided that:\nthe sum of all x\"\"s is from 2 to about 200, preferably from about 3 to about 150, more preferably from about 5 to about 120, and most preferably from about 5 to about 100;\nany basic amine site on the polymer, oligomer, or copolymer may be optionally protonated, alkylated, or quaternized with groups selected from the group consisting of H, CH3, alkyl, hydroxyalkyl, benzyl and mixtures thereof;\nany amine site may be optionally alkoxylated; and\nwhen two R1 groups are attached to a common nitrogen the two R1s may form a cyclic structure selected from the group consisting of C5-C8 alkylene, and C4-C7 alkyleneoxyalkylene.\nPreferred amino acid based polymer, oligomer or copolymer materials for use herein include lysine and hexamethylenediamine condensates of the general formula: \nwherein x and y are individually from 0 to about 50, provided that x+y greater than 0. As is illustrated in Example I below, the lysine and hexamethylenediamine condensates of the present invention are often a mixture of various molecules having different values for x+y. Especially preferred lysine and hexamethylenediamine condensates have x+y equal to from about 2 to about 5. The above, simplified structure depiction is meant to include also those oligomers and polymers with more branching which arise from amide bond linkages formed at the alpha aimino group and from amide bond linkages formed at both amino groups of a single lysine-derived unit.\nCompositions which are also preferred for use herein are lysine and octylamine condensates of the general formula: \nwherein x is from about 1 to about 50.\nAdditionally preferred compositions for use herein are lysine and tryptophan condensates of the general formula: \nWherein p and q independently are from about 1 to about 50.\nSynthesis of the amino acid based polymer, oligomer or copolymer materials defined herein will be apparent to those skilled in the art in light of Example I wherein the synthesis of 5:1, L-Lysine: 1,6-Hexanediamine is exemplified. Example I details the condensation reaction of lysine and hexanediamine which produces a preferred product for achieving the desired fabric appearance and integrity benefits of the present invention. However, other amines are equally suitable for condensation reactions with amino acids such as lysine, and those suitable amines which are readily available and/or possess interesting structural characteristics include:\naliphatic and cycloaliphatic amines, preferably methylamine, ethylamine, proplyamine, butylamine, pentylamine, hexylamine, heptylamine, ethylhexylamine, octylamine, nonylamine, decylamine, undecylamine, dodecylamine, tridecylamine, stearylamine, palmitylamine, 2-ethylhexylamine, isononylamine, hexamethyleneimine, dimethylamine, diethylamine, dipropylamine, dibutylamine, dihexylamine, ditridecylamine, N-methylbutylamine, N-ethylbutylamine;\nalkoxyalkylamines, preferably 2-methoxyethylamine, 2-ethoxyethylamine, 3-methoxypropylamine, 3-ethoxypropylamine, 3-[(2-ethylhexyl)oxy]-1-propanamine, 3-(2-methoxyethoxy)-1-propanamine, 2-methoxy-N-(2-methoxyethyl)ethanamine;\nalicyclic amines, preferably cyclopentylamine, cyclohexylamine, N-methycyclohexylamine, N-ethylcyclohexylamine, dicyclohexylamine;\ndiamines, triamines and tetramines, preferably ethylenediamine, propylenediamine, butylenediamine, neopentyldiamine, hexamethylenediamine, octamethylenediamine, imidazole, 5-amino-1,3,-trimethylcyclohexylmethylamine, 4,4xe2x80x2-methylenebiscyclohexylamine, 4,4xe2x80x2-methylenebis(2-methylcyclohexylamine), 4,7-dioxadecyl-1,10-diamine, 4,9-dioxadodecyl-1,12-diamine, 4,7,10-trioxatridecyl-1,13-diamine, 2-(ethylamino)ethylamine, 3-(methylamino)propylamine, 3-(cyclohexylamino)propylamine, 3-(2-aminoethyl) aminopropylamine, 2-(diethylamino)ethylamine, 3-(dimethylamino)propylamine;\ndimethyldipropylenetriamine, 4-aminomethyloctane-1,8-diamine, 3-(diethylamino)propylamine, N,N-diethyl-1,4-pentanediamine, diethylenetriamine, dipropylenetriamine, bis(hexamethylene)triamine, aminoethylpiperazine, aminopropylpiperazine, N,N-bis(aminopropyl)methylamine, N,N-bis(aminopropyl)ethylamine, N,N-bis(aminoproply)methylamine, N,N-bis(aminopropyl)ethylamine, N,N-bis(aminopropyl)hexylamine, N,N-bis(aminopropyl)octylamine, N,N-dimethyldipropylenetriamine, N,N-bis(3-dimethylaminopropyl)amine, N,Nxe2x80x3-1,2-ethanediylbis-(1,3-propanediamine), N-(aminoethyl)piperazine, N-(2-imidazole)piperazine, N-ethylpiperazine, N-(hydroxyethyl)piperazine, N-(aminoethyl)piperazine, N-(aminopropyl)piperazine, N-(aminoethyl)morpholine, N-(aminopropyl)morpholine, N-(aminoethyl)imidazole, N-(aminopropyl)imidazole, N-(aminoethyl)hexamethylenediamine, N-(aminopropyl)hexamethylenediamine, N-(aminoethyl)ethylenediamine, N-(aminopropyl)ethylenediamine, N-(aminoethyl)butylenediamine, N-(aminopropyl)butylenediamine, bis(aminoethyl)piperazine, bis(aminopropyl)piperazine, bis(aminoethyl)hexamethylenediamine, bis(aminopropyl)hexamethylenediamine, bis(aminoethyl)ethylenediamine, bis(aminopropyl)ethylenediamine, bis(aminoethyl)butylenediamine, bis(aminopropyl)butylenediamine;\naliphatic amino alcohols, preferably 2-aminoethanol, 3-amino-1propanol, 1-amino-2-propanol, 2-(2-aminoethoxy)ethanol, 2-[(2-aminoethyl)amino]ethanol, 2-methylaminoethanol, 2-(ethylamino)ethanol, 2-butylaminoethanol, diethanolamine, 3-[(2-hydroxyethyl)amino]-1-propanol, diisopropanolamine, bis(hydroxyethyl)aminoethylamine, bis(hydroxypropyl)aminoethylamine, bis(hydroxyethyl)aminopropylamine, bis(hydroxypropyl)aminopropylamine;\nmonoamino carboxylic acids, preferably glycine, alanine, sarcosine, asparagine, glutamine, 6-aminocaproic acid, caprolactam, 4-aminobutyric acid; and\nglucosamine, melamine, urea, guanidine, polyguanides, piperidine, morpholine, 2,6-dimethylmorpholine.\nB) Detersive Surfactant\nThe detergent compositions herein comprise from about 1% to 80% by weight of a detersive surfactant. Preferably such compositions comprise from about 5% to 50% by weight of surfactant. Detersive surfactants utilized can be of the anionic, nonionic, zwitterionic, amnpholytic or cationic type or can comprise compatible mixtures of these types. Detergent surfactants useful herein are described in U.S. Pat. No. 3,664,961, Norris, issued May 23, 1972, U.S. Pat. No. 3,919,678, Laughlin et al., issued Dec. 30, 1975, U.S. Pat. No. 4,222,905, Cockrell, issued Sep. 16, 1980, and in U.S. Pat. No. 4,239,659, Murphy, issued Dec. 16, 1980. All of these patents are incorporated herein by reference. Of all the surfactants, anionics and nonionics are preferred.\nUseful anionic surfactants can themselves be of several different types. For example, water-soluble salts of the higher fatty acids, i.e., xe2x80x9csoapsxe2x80x9d, are useful anionic surfactants in the compositions herein. This includes alkali metal soaps such as the sodium, potassium, ammonium, and alkylolammonium salts of higher fatty acids containing from about 8 to about 24 carbon atoms, and preferably from about 12 to about 18 carbon atoms. Soaps can be made by direct saponification of fats and oils or by the neutralization of free fatty acids. Particularly useful are the sodium and potassium salts of the mixtures of fatty acids derived from coconut oil and tallow, i.e., sodium or potassium tallow and coconut soap.\nAdditional non-soap anionic surfactants which are suitable for use herein include the water-soluble salts, preferably the alkali metal, and ammonium salts, of organic sulfuric reaction products having in their molecular structure an alkyl group containing from about 10 to about 20 carbon atoms and a sulfonic acid or sulfuric acid ester group. (Included in the term xe2x80x9calkylxe2x80x9d is the alkyl portion of acyl groups.) Examples of this group of synthetic surfactants are a) the sodium, potassium and ammonium alkyl sulfates, especially those obtained by sulfating the higher alcohols (C8-C18 carbon atoms) such as those produced by reducing the glycerides of tallow or coconut oil; b) the sodium, potassium and ammonium alkyl polyethoxylate sulfates, particularly those in which the alkyl group contains from 10 to 22, preferably from 12 to 18 carbon atoms, and wherein the polyethoxylate chain contains from 1 to 15, preferably 1 to 6 ethoxylate moieties; and c) the sodium and potassium alkylbenzene sulfonates in which the alkyl group contains from about 9 to about 15 carbon atoms, in straight chain or branched chain configuration, e.g., those of the type described in U.S. Pat. Nos. 2,220,099 and 2,477,383. Especially valuable are linear straight chain alkylbenzene sulfonates in which the average number of carbon atoms in the alkyl group is from about 11 to 13, abbreviated as C11-13 LAS.\nPreferred nonionic surfactants are those of the formula R1(OC2H4)nOH, wherein R1 is a C10-C16 alkyl group or a C8-C12 alkyl phenyl group, and n is from 3 to about 80. Particularly preferred are condensation products of C12-C15 alcohols with from about 5 to about 20 moles of ethylene oxide per mole of alcohol, e.g., C12-C13 alcohol condensed with about 6.5 moles of ethylene oxide per mole of alcohol.\nAdditional suitable nonionic surfactants include polyhydroxy fatty acid amides of the formula: \nwherein R is a C9-17 alkyl or alkenyl, R1 is a methyl group and Z is glycityl derived from a reduced sugar or alkoxylated derivative thereof. Examples are N-methyl N-1-deoxyglucityl cocoamide and N-methyl N-1-deoxyglucityl oleamide. Processes for making polyhydroxy fatty acid amides are known and can be found in Wilson, U.S. Pat. No. 2,965,576 and Schwartz, U.S. Pat. No. 2,703,798, the disclosures of which are incorporated herein by reference.\nPreferred surfactants for use in the detergent compositions described herein are amine based surfactants of the general formula: \nwherein R1 is a C6-C12 alkyl group; n is from about 2 to about 4, X is a bridging group which is selected from NH, CONH, COO, or O or X can be absent; and R3 and R4 are individually selected from H, C1-C4 alkyl, or (CH2xe2x80x94CH2xe2x80x94O(R5)) wherein R5 is H or methyl. Especially preferred amines based surfactants include the following:\nR1xe2x80x94(CH2)2-NH2\nR1xe2x80x94Oxe2x80x94(CH2)3-NH2\nxe2x80x83R1xe2x80x94C(O)xe2x80x94NHxe2x80x94(CH2)3-N(CH3)2 \nwherein R1 is a C6-C12 alkyl group and R5 is H or CH3. Particularly preferred amines for use in the surfactants defined above include those selected from the group consisting of octyl amine, hexyl amine, decyl amine, dodecyl amine, C8-C12 bis(hydroxyethyl)amine, C8-C12 bis(hydroxyisopropyl)amine, C8-C12 amido-propyl dimethyl amine, or mixtures thereof.\nIn a highly preferred embodiment, the amine based surfactant is described by the formula:\nR1xe2x80x94C(O)xe2x80x94NHxe2x80x94(CH2)3xe2x80x94N(CH3)2\nwherein R1 is C8-C12 alkyl.\nC) Detergent Builder\nThe detergent compositions herein may also comprise from about 0.1% to 80% by weight of a detergent builder. Preferably such compositions in liquid form will comprise from about 1% to 10% by weight of the builder component. Preferably such compositions in granular form will comprise from about 1% to 50% by weight of the builder component. Detergent builders are well known in the art and can comprise, for example, phosphate salts as well as various organic and inorganic nonphosphorus builders.\nWater-soluble, nonphosphorus organic builders useful herein include the various alkali metal, ammonium and substituted ammonium polyacetates, carboxylates, polycarboxylates and polyhydroxy sulfonates. Examples of polyacetate and polycarboxylate builders are the sodium, potassium, lithium, ammonium and substituted ammonium salts of ethylene diamine tetraacetic acid, nitrilotriacetic acid, oxydisuccinic acid, mellitic acid, benzene polycarboxylic acids, and citric acid. Other suitable polycarboxylates for use herein are the polyacetal carboxylates described in U.S. Pat. No. 4,144,226, issued Mar. 13, 1979 to Crutchfield et al., and U.S. Pat. No. 4,246,495, issued Mar. 27, 1979 to Crutchfield et al., both of which are incorporated herein by reference. Particularly preferred polycarboxylate builders are the oxydisuccinates and the ether carboxylate builder compositions comprising a combination of tartrate monosuccinate and tartrate disuccinate described in U.S. Pat. No. 4,663,071, Bush et al., issued May 5, 1987, the disclosure of which is incorporated herein by reference.\nExamples of suitable nonphosphorus, inorganic builders include the silicates, aluminosilicates, borates and carbonates. Particularly preferred are sodium and potassium carbonate, bicarbonate, sesquicarbonate, tetraborate decahydrate, and silicates having a weight ratio of SiO2 to alkali metal oxide of from about 0.5 to about 4.0, preferably from about 1.0 to about 2.4. Also preferred are aluminosilicates including zeolites. Such materials and their use as detergent builders are more fully discussed in Corkill et al., U.S. Pat. No. 4,605,509, the disclosure of which is incorporated herein by reference. Also discussed in U.S. Pat. No. 4,605,509 are crystalline layered silicates which are suitable for use in the detergent compositions of this invention.\nD) Optional Detergent Ingredients\nIn addition to the surfactants, builders and amino acid based polymer, oligomer or copolymer materials hereinbefore described, the detergent compositions of the present invention can also include any number of additional optional ingredients. These include conventional detergent composition components such as enzymes and enzyme stabilizing agents, suds boosters or suds suppressers, anti-tarnish and anticorrosion agents, soil suspending agents, soil release agents, germicides, pH adjusting agents, non-builder alkalinity sources, chelating agents, organic and inorganic fillers, solvents, hydrotropes, optical brighteners, dyes and perfumes.\npH adjusting agents may be necessary in certain applications where the pH of the wash solution is greater than about 10.0 because the fabric integrity benefits of the defined compositions begin to diminish at a higher pH. Hence, if the wash solution is greater than about 10.0 after the addition of the amino acid based polymer, oligomer or copolymer materials of the present invention a pH adjuster should be used to reduce the pH of the washing solution to below about 10.0, preferably to a pH of below about 9.5 and most preferably below about 7.5. Suitable pH adjusters will be known to those skilled in the art.\nNormally, a preferred optional ingredient for incorporation into detergent compositions is a bleaching agent, e.g., a peroxygen bleach. However, many common bleaching agents will degrade some, but not all, of the amino acid based fabric treatment materials of the present invention. Hence, before adding a bleaching agent to a detergent composition comprising an amino acid based fabric treatment material as defined herein compatibility between the bleaching agent and the amino acid based fabric treatment material must be investigated.\nAnother highly preferred optional ingredient in the detergent compositions herein is a detersive enzyme component. While it is known that some enzymes will degrade the peptide bonds of amino acids, the amino acid based polymer, oligomer or copolymer materials defined herein do not exhibit such degradation in the presence of enzymes. Hence, enzymes can be added to detergent compositions which comprise the amino acid based fabric treatment materials of the present invention with substantially no degradation.\nEnzymes can be included in the present detergent compositions for a variety of purposes, including removal of protein-based, carbohydrate-based, or triglyceride-based stains from substrates, for the prevention of refugee dye transfer in fabric laundering, and for fabric restoration. Suitable enzymes include proteases, amylases, lipases, cellulases, peroxidases, and mixtures thereof of any suitable origin, such as vegetable, animal, bacterial, fungal and yeast origin. Preferred selections are influenced by factors such as pH-activity and/or stability, optimal thermostability, and stability to active detergents, builders and the like. In this respect bacterial or fungal enzymes are preferred, such as bacterial amylases and proteases, and fungal cellulases. xe2x80x9cDetersive enzymexe2x80x9d, as used herein, means any enzyme having a cleaning, stain removing or otherwise beneficial effect in a laundry detergent composition. Preferred enzymes for laundry purposes include, but are not limited to, proteases, cellulases, lipases, amylases and peroxidases.\nEnzymes are normally incorporated into detergent compositions at levels sufficient to provide a xe2x80x9ccleaning-effective amountxe2x80x9d. The term xe2x80x9ccleaning-effective amountxe2x80x9d refers to any amount capable of producing a cleaning, stain removal, soil removal, whitening, deodorizing, or freshness improving effect on substrates such as fabrics. In practical terms for current commercial preparations, typical amounts are up to about 5 mg by weight, more typically 0.01 mg to 3 mg, of active enzyme per gram of the detergent composition. Stated otherwise, the compositions herein will typically comprise from 0.001% to 5%, preferably 0.01%-1% by weight of a commercial enzyme preparation. Protease enzymes are usually present in such commercial preparations at levels sufficient to provide from 0.005 to 0.1 Anson units (AU) of activity per gram of composition. Higher active levels may be desirable in highly concentrated detergent formulations.\nSuitable examples of proteases are the subtilisins which are obtained from particular strains of B. subtilis and B. licheniformis. One suitable protease is obtained from a strain of Bacillus, having maximum activity throughout the pH range of 8-12, developed and sold as ESPERASE(copyright) by Novo Industries A/S of Denmark, hereinafter xe2x80x9cNovoxe2x80x9d. The preparation of this enzyme and analogous enzymes is described in GB 1,243,784 to Novo. Other suitable proteases include ALCALASE(copyright) and SAVINASE(copyright) from Novo and MAXATASE(copyright) from International Bio-Synthetics, Inc., The Netherlands; as well as Protease A as disclosed in EP 130,756 A, Jan. 9, 1985 and Protease B as disclosed in EP 303,761 A, Apr. 28, 1987 and EP 130,756 A, Jan. 9, 1985. See also a high pH protease from Bacillus sp. NCIMB 40338 described in WO 9318140 A to Novo. Enzymatic detergents comprising protease, one or more other enzymes, and a reversible protease inhibitor are described in WO 9203529 A to Novo. Other preferred proteases include those of WO 9510591 A to Procter and Gamble. When desired, a protease having decreased adsorption and increased hydrolysis is available as described in WO 9507791 to Procter and Gamble. A recombinant trypsin-like protease for detergents suitable herein is described in WO 9425583 to Novo.\nCellulases usable herein include both bacterial and fungal types, preferably having a pH optimum between 5 and 10. U.S. Pat. No. 4,435,307, Barbesgoard et al., Mar. 6, 1984, discloses suitable fungal cellulases from Humicola insolens or Humicola strain DSM1800 or a cellulase 212-producing fungus belonging to the genus Aeromonas, and cellulase extracted from the hepatopancreas of a marine mollusk, Dolabella Auricula Solander. Suitable cellulases are also disclosed in GB-A-2.075.028; GB-A-2.095.275 and DE-OS-2.247.832. CAREZYME(copyright) and CELLUZYME(copyright) (Novo) are especially useful. See also WO 9117243 to Novo.\nSuitable lipase enzymes for detergent usage include those produced by microorganisms of the Pseudomonas group, such as Pseudomonas stutzeri ATCC 19.154, as disclosed in GB 1,372,034. See also, the lipase in Japanese Patent Application 53,20487, laid open Feb. 24, 1978. This lipase is available from Amano Pharmaceutical Co. Ltd., Nagoya, Japan, under the trade name Lipase P xe2x80x9cAmano,xe2x80x9d or xe2x80x9cAmano-P.xe2x80x9d Other suitable commercial lipases include Amano-CES, lipases ex Chromobacter viscosum, e.g. Chromobacter viscosum var. lipolyticum NRRLB 3673 from Toyo Jozo Co., Tagata, Japan; Chromobacter viscosum lipases from U.S. Biochemical Corp., U.S.A. and Disoynth Co., The Netherlands, and lipases ex Pseudomonas gladioli. LIPOLASE(copyright) enzyme derived from Humicola lanuginosa and commercially available from Novo, see also EP 341,947, is a preferred lipase for use herein.\nThe enzyme-containing compositions herein may optionally also comprise from about 0.001% to about 10%, preferably from about 0.005% to about 8%, most preferably from about 0.01% to about 6%, by weight of an enzyme stabilizing system. The enzyme stabilizing system can be any stabilizing system which is compatible with the detersive enzyme. Such a system may be inherently provided by other formulation actives, or be added separately, e.g., by the formulator or by a manufacturer of detergent-ready enzymes. Such stabilizing systems can, for example, comprise calcium ion, boric acid, propylene glycol, short chain carboxylic acids, boronic acids, and mixtures thereof, and are designed to address different stabilization problems depending on the type and physical form of the detergent composition.\nE) Detergent Composition Preparation\nThe detergent compositions according to the present invention can be in liquid, paste or granular form. Such compositions can be prepared by combining the essential and optional components in the requisite concentrations in any suitable order and by any conventional means.\nGranular compositions, for example, are generally made by combining base granule ingredients, e.g., surfactants, builders, water, etc., as a slurry, and spray drying the resulting slurry to a low level of residual moisture (5-12%). The remaining dry ingredients, e.g., granules of the essential amino acid based fabric treatment materials, can be admixed in granular powder form with the spray dried granules in a rotary mixing drum. The liquid ingredients, e.g., solutions of the essential amino acid based fabric treatment materials, enzymes, binders and perfumes, can be sprayed onto the resulting granules to form the finished detergent composition. Granular compositions according to the present invention can also be in xe2x80x9ccompact formxe2x80x9d, i.e. they may have a relatively higher density than conventional granular detergents, i.e. from 550 to 950 g/l. In such case, the granular detergent compositions according to the present invention will contain a lower amount of xe2x80x9cinorganic filler saltxe2x80x9d, compared to conventional granular detergents; typical filler salts are alkaline earth metal salts of sulphates and chlorides, typically sodium sulphate; xe2x80x9ccompactxe2x80x9d detergents typically comprise not more than 10% filler salt.\nLiquid detergent compositions can be prepared by admixing the essential and optional ingredients thereof in any desired order to provide compositions containing components in the requisite concentrations. Liquid compositions according to the present invention can also be in xe2x80x9ccompact formxe2x80x9d, in such case, the liquid detergent compositions according to the present invention will contain a lower amount of water, compared to conventional liquid detergents. Addition of the amino acid based polymer, oligomer or copolymer materials to liquid detergent or other aqueous compositions of this invention may be accomplished by simply mixing into the liquid solutions the desired amino acid based fabric treatment materials.\nF) Fabric Laundering Method\nThe present invention also provides a method for laundering fabrics in a manner which imparts fabric appearance benefits provided by the amino acid based polymer, oligomer or copolymer materials used herein. Such a method employs contacting these fabrics with an aqueous washing solution formed from an effective amount of the detergent compositions hereinbefore described or formed from the individual components of such compositions. Contacting of fabrics with washing solution will generally occur under conditions of agitation although the compositions of the present invention may also be used to form aqueous unagitated soaking solutions for fabric cleaning and treatment. As discussed above, it is preferred that the washing solution have a pH of less than about 10.0, preferably it has a pH of about 9.5 and most preferably it has a pH of about 7.5.\nAgitation is preferably provided in a washing machine for good cleaning. Washing is preferably followed by drying the wet fabric in a conventional clothes dryer. An effective amount of a high density liquid or granular detergent composition in the aqueous wash solution in the washing machine is preferably from about 500 to about 7000 ppm, more preferably from about 1000 to about 3000 ppm.\nG) Fabric Conditioning\nThe amino acid based polymer, oligomer or copolymer materials hereinbefore described as components of the laundry detergent compositions herein may also be used to treat and condition fabrics and textiles in the absence of the surfactant and builder components of the detergent composition embodiments of this invention. Thus, for example, a fabric conditioning composition comprising only the amino acid based fabric treatment materials themselves, or comprising an aqueous solution of the amino acid based fabric treatment materials, may be added during the rinse cycle of a conventional home laundering operation in order to impart the desired fabric appearance and integrity benefits hereinbefore described."}
{"text": "1. Field of the Invention\nThis invention relates to a polishing and scrubbing apparatus similar to polishing pads used with electric or propane fueled internal combustion engine polishing machines.\n2. Statement of the Prior Art\nThe prior art shows numerous structures of polishing pads for use with floor polishing machines. These pads conventionally have a central aperture for attaching them to the undersurface of a polishing machine. One of the disadvantages of these conventional pads is that they do not present cutting surfaces to cut wax from the floor. An additional disadvantage of conventional pads is that they create a great amount of drag on the machine often resulting in overheating and consequent overloading of electrical circuits.\nThe present invention overcomes the disadvantages in conventional pads by providing a polishing apparatus having multiple cutting edges to remove hardened wax on a floor as well as reduced surface area thus preventing drag.\nRepresentative of prior art devices are disclosed in the following U.S. Pat. Nos. and copies are furnished for the record: 3,183,542 on 5/1965 to Anders; 3,793,665 on 2/1974 to Thielen; 4,307,480 on 12/1891 to Fallen; 4,502,174 on 5/1985 to Rones."}
{"text": "The mobile telephone industry has been associated with tremendous growth over the last several years. For instance, in the recent past, mobile telephones were costly, bulky and awkward to transport. Moreover, network coverage was not extensive enough to enable robust service. Only areas associated with dense population were provided with extensive wireless network coverage. In contrast, today's mobile phones (and other portable devices) have decreased in both size and cost and can be utilized as full-service computing machines. For example, many of the most recent and advanced mobile phones can be associated with word processing software, accounting software, and various other types of software. Network coverage has expanded to cover millions, if not billions, of users.\nAdvances in technology relating to mobile devices in general, and mobile phones in particular, continue to occur. For example, recently mobile telephones have been designed to communicate over disparate networks. For example, a dual mode handset can connect to a cellular network to effectuate communications between a user of the mobile phone and another phone device, and can further connect via WiFi to a wireless local access network (LAN) and thereafter utilize the Voice over Internet Protocol (VoIP) to effectuate communication between users. Use of VoIP is often desirable to users as it is associated with lower costs than employing a cellular network. In fact, some users may consider phone calls made over VoIP (or other IP-based network) completely free, despite the fact that they pay for Internet service.\nImplementation of this dual mode service (DMS) is due at least in part to the Third Generation Partnership Project (3GPP), which have created specifications that define a mechanism that provides signal integrity for session initial protocol (SIP) signals between an IP multimedia subsystem (IMS) (P-SCCF) and user equipment (UE) (e.g., a mobile phone, a personal digital assistant, . . . ). This integrity prevents identity spoofing, man-in-the-middle attacks, and the like. The IMS represents a 3GPP and 3GPP2 effort to define an all-IP-based wireless network as a replacement for the various voice, data, signaling, and control network elements currently in existence. Furthermore, the IMS enables support for IP multimedia applications within the Universal Mobile Telecommunications System (UMTS). The UMTS is a 3G broadband packet-based transmission of text, digitized voice, video, and multimedia that offers a consistent set of services to mobile computer and phone users regardless of their physical location.\nThe telecom industry is currently shifting towards all IP-systems, thereby rendering dual mode service handsets an important tool (as they are compatible with existing cellular systems and emerging IP-systems). This shift is driven by desires to reduce costs and create new streams of revenue while protecting an operator business model. IMS is a new service domain that facilitates this shift by enabling convergence of data, speech, and network technology over an IP-based infrastructure. For users, IMS-based services enable transmittal and receipt of various data at significantly reduced cost, including voice, text, pictures, video, and/or any combination thereof in a highly personalized and secure manner. In summary, IMS is designed to bridge the gap between existing, traditional telecommunications technology and Internet technology that increased bandwidth does not provide.\nAs stated above, these emerging IP-based technologies have created demand for dual mode services, and thus for dual mode handsets. Using this technology, subscribers can employ WiFi to effectuate voice calls, transmission of data, and the like. In more detail, a subscriber can connect to a LAN by way of WiFi. Upon such connection, subscribers can employ services offered by their service provider.\nAs VoIP systems and dual mode handsets have become more accepted and prevalent, these systems have been required to provide the same capabilities as traditional telephone networks, including emergency services. During an emergency situation, a user can be flustered or injured and unable to provide accurate location information, whether due to age, infirmity, injury or ignorance. Typically, traditional fixed line telephone networks can connect a caller to local emergency services and automatically provide location information to the emergency service center. In the past, VoIP systems have not reliably provided location information or even connection to emergency services such as Enhanced 911 (E911). However, government regulations adopted by the Federal Communications Commission (FCC) now require VoIP providers to supply E911 service capabilities to their subscribers."}
{"text": "Over the last decade, there has been a substantial increase in the use and deployment of wireless client devices, from dual-mode smartphones to tablets capable of operating in accordance with a particular Institute of Electrical and Electronics Engineers (IEEE) standard. With “wireless” becoming the de-facto medium for connectivity among users, it has become increasingly important for network systems to intelligently manage connections.\nIn some systems, multiple access points may be used to cover a particular area. Each of these access points may operate using separate wireless channels. However, in these traditional systems assignment of channels is performed based on the assumption that each channel is equal. This assumption may result in a network system that fails to adequately support particular client devices and/or provides poor network coverage.\nThe approaches described in this section are approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section."}
{"text": "The instant invention relates to exercise apparatus for exercising the adductor muscles of the legs and more particularly to an adductor exercise apparatus having means for adjusting the angular starting position of the leg receiving members.\nExercise machines, and more particularly, adductor exercise machines, have heretofore been known in the art. In this connection, the U.S. Patents to Scott No. 4,022,463; DeNiro No. 4,892,304; Dela Rosa No. 4,877,239; and Goodman No. 5,026,049 represent the closest prior art to the subject invention of which the applicant is aware.\nThe patent to Scott discloses an exercise apparatus for exercising the adductor muscles of the legs, however it does not disclose any means for adjusting the starting position of the leg receiving means. The patent to DeNiro discloses an exercise apparatus including a plurality of movable parts which are adjusted by means of pin and hole mechanisms. The patent to Dela Rosa discloses a stretching apparatus having an elongated threaded adjusting shaft in the form of a worm gear for adjusting the degree of split of the leg receiving elements. The patent to Goodman shows leg receiving pads which are adjustable to accommodate different size persons.\nThe instant invention provides an adductor exercise apparatus which includes an adjustment mechanism for adjusting the angular starting positions of the leg receiving assemblies. More specifically, the adductor exercise apparatus includes a base, two support legs which extend outwardly from opposite sides of the base, a seat on the base, and two leg receiving assemblies for receiving the legs of a user seated on the seat. The leg assemblies are pivotably mounted to the base so that they are pivotably movable between a spread apart position and a parallel, together position. The leg assemblies include a pad assembly including two pads which respectively engage with the thigh and calf portions of the user's leg. The calf pad is slidably movable with respect to the thigh pad. The apparatus further includes a cabled weight assembly for normally urging the leg assemblies toward the spread apart position. The adjustment mechanism consists of a pair of movable adjustment arms each having a first end which engages the corresponding leg receiving assembly, and a second end which is slidably received and secured in a sleeve assembly mounted on the support leg adjacent to the respective leg receiving assembly. The arms are slidably adjustable in the sleeve assemblies to a plurality of predetermined longitudinal positions. In this manner, as the leg receiving assemblies are urged toward their normal spread apart position, the leg receiving assemblies engage with the projecting ends of the adjustment arms to position the leg receiving assemblies in corresponding angular starting positions. The cabled weight assembly includes stacked weight members and a cable system which extends around a series of pulleys for translating pivoting movement of the leg receiving assemblies into vertical movement of the weight members. The weight members are lifted by an elongated weight bar which is attached to the cable system.\nAccordingly, it is an object of the instant invention to provide an exercise apparatus for exercising the adductor muscles of the leg.\nIt is another object to provide an adductor exercise apparatus having an adjustment mechanism for adjusting the angular starting positions of the leg receiving assemblies.\nIt is yet another object to provide an adductor exercise apparatus including leg receiving means having two separate pads for engaging with the thigh and calf portions of the user's leg.\nIt is still another object to provide an adductor exercise apparatus having a cabled weight assembly and an elongated weight bar for lifting the weight members of the weight assembly.\nOther objects, features and advantages of the invention shall become apparent as the description thereof proceeds when considered in connection with the accompanying illustrative drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a press for the stamping of plates, particularly vehicle license plates having license numbers. The license plates are made from sheet metal material, itself in the form of strip material or individual plates, that is fed in synchronized manner into the press. The press includes a press stand comprised of a press table and a press stamp, as well as interchangeable stamping tools.\n2. Description of the Related Art\nPresses of the above type, as found in DE 32 03 801 C2, are furnished with tools, embodied as block tools, having a uniform width. The tools are used for stamping license numbers (having letters and numbers) in vehicle license plates and do so in registered layout, or in a registered print style, as known in Germany.\nIn other countries such as Austria, the figures on the license plate are stamped in a non-registered layout, or in a free print style. In other words, the stamping tools have a width that can be tailored to the width of letters, numbers and, if need be, coats of arms.\nThe present invention seeks to develop a press for the stamping of plates, particularly vehicle license plates with license numbers in any linear order, so that it can be furnished with stamping tools of uniform width (i.e. uniform registration) as well as with tools having different widths."}
{"text": "Insulin-like growth factor-1 belongs to a heterogeneous family of peptides which share some of the biological and chemical properties of insulin, but which are antigenically distinct from insulin. Currently available experimental evidence suggests that IGF-1 promotes growth by mediating the effects of growth hormone. Thus, such processes as skeletal growth, call replication and other growth related processes are affected by IGF-1 levels. Physiological concentrations of IGF-1 have been shown to be influenced by such conditions as thyroid disease, diabetes and malnutrition (see Preece (1983) Horm. Blood, 4: 108). IGF-1 has also been shown to act synergistically with other growth factors, for example, in accelerating the healing of soft and mesenchymal tissue wounds (see Lynch et al. (1989) J. Clin., Periodontol, 16: 545; and Lynch et al. (1987) Proc. Nat. Acad. Sci. USA, 84: 7696), and in enhancing the growth of mammalian cells in serum-free tissue culture medium (see Burleigh et al. (1986) American Biotech. Lab., 4: 48).\nConsidering the many clinical and research applications of IGF-1, a ready supply of IGF-1 would be of great value to the medical and biotechnology fields. Since isolation from natural sources is technically difficult, expensive, and time consuming, recent efforts have centered on the development of efficient recombinant methods for the production of IGF-1.\nAmong host cells that have been used for the production of heterologous proteins, E. coli and Saccharomyces cerevisiae (Baker's yeast) are probably the best characterized. Insulin-like growth factor-1 (IGF-1), which is a polypeptide of 70 amino acids with a molecular weight of 7648 daltons, is a single chain protein that has three intrachain disulfide bridges. These disulfide bonds, along with numerous hydrogen bonds and hydrophilic interactions, maintain the compact tertiary structure of this molecule. E. coli, however, does not possess the ability to produce disulfide bonds in proteins, so that proteins, such as IGF-1, that include disulfide bonds, when cloned into and expressed in E. coli, frequently are not stable and tend to aggregate into inactive complexes. In addition, IGF-1 produced in E. coli has to be extracted and treated with oxidizing agents to produce the disulfide bonds. Upon reduction and reoxidation, IGF-1 refolds in a variety of ways, forming as many as 15 monomeric configurations (Meng et al. (1988) J. Chrom., 443: 183) because cell breakage and too rapid formation of disulfide linkages results in random disulfide bond formation. In order to produce biologically active IGF-1, the resulting mixture of 15 different forms of IGF-1 must be separated. Consequently, the yield of purified product is very low (Grossgian (1985) Gene, 18: 199).\nFurthermore, since E. coli is a prokaryote, in order to produce IGF-1 molecules which contain the authentic N-terminal glycine, and not the initiating methionine present on the primary translation product, it is necessary to express IGF-1 in E. coli as a fusion protein. Cleavage of mature IGF-1 from the initially produced fusion protein necessitates an additional step in the production process. Consequently, attempts to produce this peptide by recombinant means in E. coli host expression systems results in a complex mixture of product forms which must be separated for further use (see, Grossgian (1985) Gene, 18: 199).\nEukaryotic host cells, including yeast cells, thus, are the host cells of choice for the expression of many eukaryotic proteins. Yeast host cells offer clear advantages over bacteria in the production of heterologous proteins, including their ability to properly process pre-pro-heterologous proteins and secrete heterologous proteins into the culture medium. Secretion of proteins from cells is often superior to production of proteins in the cytoplasm because secreted products are obtained in a higher degree of initial purity and further purification of the secreted products is made easier by the absence of cellular debris. In addition, the secretory pathway of the cell and the extracellular medium tend to be oxidizing environments which support disulfide bond formation necessary for proper folding of many proteins (Smith, et al, (1985) Science 229: 1219); whereas, the cytoplasm is a reducing environment in which disulfide bonds do not form. Thus, for production of sulfhydryl-rich proteins that rely on disulfide bonds to maintain the correct tertiary structures, there is a compelling need to develop eukaryotic hosts capable of secreting such proteins into the culture medium. Therefore, production of sulfhydryl-rich proteins, such as IGF-1, that contain appropriately formed disulfide bonds, can best be achieved by transit through the secretory pathway.\nIGF-1 has been cloned into and expressed using S. cerevisiae host cells by introducing DNA encoding IGF-1 on autonomously replicating extrachromosomal elements. Gellerfors et al. ((1989) J. Biol, Chem., 264: 11444-11449) describes the production of IGF-1 in S. cerevisiae under the control of the S. cerevisiae actin promoter. The IGF-1 product is encoded by autonomously replicating plasmid-borne DNA. In a similar study, Bayne et al. ((1988) Gene 66: 235-244) describes the production of IGF-1 in S. cerevisiae under the control of the S. cerevisiae alpha mating factor promoter. The latter, however, reports yields of IGF-1 of only about 2 mg of IGF-1 per liter of fermentation broth.\nIn view of this low yield and the problems generally encountered with up-scaling the production of heterologous proteins in autonomous plasmid-based yeast systems, such as loss of selection for plasmid maintenance and problems concerning plasmid distribution, copy number and stability in fermentors operated at high cell density, there is a need to develop more efficient means for producing large quantities of biologically active IGF-1.\nTherefore, it is an object of this invention to provide host cells and expression vectors that stably express IGF-1 and that secrete high concentrations of biologically active IGF-1.\nIt is another object of this invention to provide an expression system for the production of biologically active IGF-1 that, not only secretes high concentrations of biologically active IGF-1, but that can be readily scaled up to produce large quantities of such IGF-1."}
{"text": "The building industry including home construction involves fiber glass material being used as insulation above the ceilings and in the walls. This insulation can be either blown in place or sheet laid. The insulation material in either case produces particles injurious to a person's skin, nose, lungs and eyes. Conventionally, protection against the fiber glass particles involves the workman wearing goggles and rubber gloves on the hands. This procedure does not work satisfactorily as the particles can contact the unprotected part of the workman's body and seep behind the goggles through small openings. The rubber gloves interfere with finger gripping of tools and building materials.\nWhat is needed is a garment that will fully protect the person's body but yet allow for complete vision and use of the hands including the fingers."}
{"text": "Cancer is one of the most widespread diseases afflicting mankind and a major cause of death worldwide. In an effort to find an effective treatment or a cure for one or more of the many different types of cancer, over the last couple of decades, numerous groups have invested a tremendous amount of time, effort and financial resources. However, to date, of the available cancer treatments and therapies, only a few offer any considerable degree of success.\nCancer is often characterized by unregulated cell proliferation. Damage to one or more genes, responsible for the cellular pathways, which control progress of proliferation through the cell cycle, typically causes the loss of normal regulation of cell proliferation. These genes code for various proteins, which participate in a cascade of events, including protein phosphorylation, leading to cell-cycling progression and cell proliferation. Various kinase proteins have been identified, which play roles in the cell cycling cascade and in protein phosphorylation in particular.\nOne class of proteins found to play a part in cell cycling and, therefore, cell proliferation is the Aurora kinase family of proteins. Aurora kinases are enzymes of the serine/threonine kinase family of proteins, which play an important role in protein phosphorylation during the mitotic phase of the cell cycle. There are three known members of the Aurora kinase family, Aurora A, Aurora B and Aurora C, also commonly referred to as Aurora 2, Aurora 1, and Aurora 3, respectively.\nThe specific function of each Aurora kinase member in mammalian cell cycle has been studied. Aurora-A is localized to the centrosome during interphase and is important for centrosome maturation and to maintain separation during spindle assembly. Aurora-B localizes to the kinetochore in the G2 phase of the cell cycle until metaphase, and relocates to the midbody after anaphase. Aurora-C was thought to function only in meiosis, but more recently has been found to be more closely related to Aurora-B, showing some overlapping functions and similar localization patterns in mitosis. Each aurora kinase appears to share a common structure, including a highly conserved catalytic domain and a very short N-terminal domain that varies in size. (See R. Giet and C. Prigent, J. Cell. Sci., 112:3591-3601 (1999)).\nAurora kinases appear to be viable targets for the treatment of cancer. Aurora kinases are overexpressed in various types of cancers, including colon, breast, lung, pancrease, prostate, bladder, head, neck, cervix, and ovarion cancers. The Aurora-A gene is part of an amplicon found in a subset of breast, colon, ovarian, liver, gastric and pancreatic tumors. Aurora-B has also been found to be overexpressed in most major tumor types. Overexpression of Aurora-B in rodent fibroblasts induces transformation, suggesting that Aurora-B is oncogenic. More recently, Aurora-B mRNA expression has been linked to chromosomal instability in human breast cancer. (Y. Miyoshi et al., Int. J. Cancer, 92:370-373 (2001)).\nFurther, inhibition of one or more of the Aurora kinases by several parties has been shown to inhibit cell proliferation and trigger apoptosis in several tumor cell lines. Particularly, inhibition of Aurora has been found to arrest cell cycling and promote programmed cell death via apoptosis. Accordingly, there has been a strong interest in finding inhibitors of Aurora kinase proteins.\nThus, the inhibition of Aurora kinases has been regarded as a promising approach for the development of novel anti-cancer agents. For example, WO 04/039774 describes aza-quinazolinones for treating cancer via inhibiton of Aurora kinase, WO 04/037814 describes indazolinones for treating cancer via inhibiton of Aurora-2 kinase, WO 04/016612 describes 2, 6, 9-substituted purine derivatives for treating cancer via inhibiton of Aurora kinase, WO 04/000833 describes tri- and tetra-substituted pyrimidine compounds useful for treating Aurora-meiated diseases, WO 04/092607 describes crystals useful for screening, designing and evaluating compounds as agonists or antagonists of Aurora kinase and U.S. Pat. No. 6,919,338 and WO 03/055491 each describe substituted quinazoline derivatives as inhibitors of Aurora-2 kinase."}
{"text": "A pivotable canopy roof of an air passenger stair or an air passenger bridge of the type mentioned in the introduction is known from DE 10 2004 016 272 B4. The design of this canopy roof has proven of great value and is often used. However, it has turned out that with strongly contoured aircraft fuselages, and here more specifically in the area of the bow side boarding access, the flexible frame with which the canopy roof rests on the outside skin of the aircraft fuselage not always rests flush with the aircraft hull on the bow side, i.e. toward the front end of the aircraft. On the one hand, this can be due to the fact that, as has already been explained, the aircraft fuselage is strongly contoured in that access area, or that more specifically with large canopy roofs, the inherent elasticity of the canopy roof and here more specifically of the articulated arms is minimal due to the stiff design of the articulated arms. This means that the deformation is not sufficient to ensure a full contact of the flexible frame with the bumper.\nIn order to always ensure a circumferential contact of the flexible frame on the outer skin of the aircraft fuselage, it is proposed according to the invention that the end member of the at least one first articulated arm is pivotable laterally via a swivel joint. Here, laterally pivotable means that the end member is pivotable in the direction of the width of the flexible frame, more specifically in the direction of the interior of the canopy roof. This means that when the flexible frame of the pivotable canopy roof comes to rest on the outer skin of the aircraft, at the moment in which the flexible frame comes in contact with the outer skin in the area of the end member, the end member of the at least one first articulated arm is pivoted slightly inward, and thus ensures full contact of the flexible frame on the outer skin of the aircraft. By pivoting the end member inward, it is ensured that the canopy roof rests on the aircraft fuselage substantially without tension.\nAdvantageous features and embodiments of the invention can be gathered from the dependent claims. In the following, the structure of the first articulated arm is described in more detail.\nThus, it is more specifically provided that in the one first articulated arm, the upper arm is connected to the lower arm by way of a hinge joint with a hinge axis that is horizontal in the mounted state, the joint axis of the pivot joint for the laterally pivotable or extendable end member running perpendicular to the hinge axis between the upper arm and the lower arm.\nAccording to another feature of the invention, the end member is connected to the lower arm by way of an intermediate member. This means that the first articulated arm consists of a total of four elements, namely an upper arm, a lower arm connected to it in an articulated manner, to which an intermediate member is connected in turn, which is preferably angled in the direction of the pivot plane of the articulated arm relative to the lower arm, and the end member that is receivable by the intermediate member so as to be laterally pivotable around a pivot joint.\nAccording to a particularly advantageous feature, the end member is pivotable out of the plane of the lower arm against the force of a spring. Thus it is ensured that when retracting the canopy roof, the end member and the first articulated arm always lie in one plane.\nFurthermore, ensuring that the end member is laterally pivotable in only one direction, namely laterally in the direction of the interior of the canopy roof is advantageous.\nThe connection between the intermediate member and the end member is configured as a laterally movable connection, the intermediate member comprising a projection that is received with a lateral clearance by the end member. This means that the projection of the intermediate member is mounted in the hollow profile of the end member with a lateral clearance, in order to permit the lateral extendibility of the end member relative to the intermediate member. An elastomer block relative to which the end member is resiliently pivotable is located between the projection and the end member configured as a hollow profile. This means that the distance between the projection of the intermediate member and the interior wall of the end member determines amongst others the degree of deflection of the end member.\nAccording to another feature of the invention, the air passenger stair or the air passenger bridge comprises an actuation device for pivoting the canopy roof, the actuation device including a pulling means at least for the one first articulated arm. A roll for the pulling means is disposed in the area of the portal frame, the pulling means being laterally articulated with the end member of the articulated arm, i.e. on the side toward which the end member does not pivot. This means that the articulation of the pulling means is eccentric. The consequence of this is that the arrangement of the pulling means on the outer side of the end member assists in the straightening of the end member.\nThe second articulated arm is also advantageously connected to the portal frame by way of a traction mechanism, e.g. a rope or a strap.\nIn particular, a roll on the portal frame for rolling up and unrolling the pulling means, i.e. for example a strap or a rope, is attributed to each pulling means, the two rolls being connected by way of a shaft to a drive, for example a tubular motor.\nIn order to initiate the extension movement of the articulated arms, the actuation device includes at least one driving member for the at least one first articulated arm, which is more specifically configured as a gas pressure spring.\nIn the following, the invention is exemplarily described in more detail based on the drawings."}
{"text": "There is such a conventionally known engine block of which cylinder block is comprised of a cylinder block body made of an aluminum alloy, and one or two or more cylinder sleeves of a cast iron formed by casting in the cylinder block body, and to which engine block a cylinder head and bearing caps for a crankshaft are assembled through a plurality of stud bolts provided on the cylinder block.\nThe conventional engine block suffers from the following problem: An explosion force applied to the cylinder head is received by the cylinder block, and the cylinder block is a member to which the cylinder head and the bearing caps are assembled. Therefore, in order to increase the strength of the cylinder block and thus the strengths of the cylinder block body and the cylinder sleeves, the thickness of them is increased. For this reason, the weight of the cylinder block is increased and hence, it is impossible to meet the demand for a reduction in weight of the engine block.\nIn a waste heat recovering device for an internal combustion engine utilizing a Rankin cycle, the following problem is encountered: If an exhaust gas from the internal combustion engine is used as a heat source, the higher the temperature of the exhaust gas is, the higher the waste heat recovering efficiency is. If each of the cylinder sleeves of the cast iron is formed by casting in the cylinder block body, as described above, the area of contact of both the members with each other is increased. For this reason, the heat of the exhaust gas is partially diffused from the cylinder sleeves through the cylinder block body, and as a result, the temperature of the exhaust gas dropps."}
{"text": "Hydraulic cylinders are mechanical actuators that get power from pressurized hydraulic fluid. A hydraulic cylinder typically includes a cylinder barrel in which a piston connected to a piston rod moves back and forth. The piston divides the hydraulic cylinder into a first chamber and a second chamber. When the hydraulic pump pushes hydraulic fluid into the first chamber, a valve in the second chamber is open allowing hydraulic fluid to drain from the second chamber into a reservoir as movement of the piston within the hydraulic cylinder increases the volume of the first chamber and correspondingly reduces the volume of the second chamber. Likewise, when the hydraulic pump pushes hydraulic fluid into the second chamber, a valve in the first chamber is open allowing hydraulic fluid to drain from the first chamber into the reservoir as movement of the piston within the hydraulic cylinder increases the volume of the second chamber and correspondingly reduces the volume of the first chamber.\nTypically, the hydraulic pump runs at a constant speed to produce hydraulic pressure. If motion is not imminent, the unused pressured hydraulic fluid is returned to the reservoir or stored in an accumulator."}
{"text": "1. Field of the Invention\nExemplary embodiments relate to instant messages, and particularly to redirecting instant messages from one system to another system.\n2. Description of Related Art\nThere exists a growing popularity in instant messaging services. Instant messaging (IM) is a form of real-time communication between two or more people based on typed text. The text is conveyed via computers connected over a network such as the Internet.\nInstant messaging offers real-time and/or near-time communication and allows easy collaboration, which might be considered more akin to genuine conversation than email's “letter” format. In contrast to e-mail, the parties know whether the peer is available because most systems allow the user to set an online status or away message so that peers are notified when the user is available, busy, or away from the computer.\nInstant messaging allows instantaneous communication between a number of parties simultaneously, by transmitting information quickly and efficiently, featuring immediate receipt of acknowledgment or reply. In certain cases IM involves additional features, which make it even more popular, i.e., to see the other party, e.g. by using web-cams, or to talk directly for free over the Internet.\nCurrently, when utilizing instant messaging from one computer and then subsequently logging into the instant messaging session from another computer, the first computer may automatically be logged off. Although new messages may be sent to the second computer, messages already delivered to the first computer would remain on the first computer screen until the user returns to the first computer and dismisses them. Thus, the delivered messages would not be delivered to the user in a timely fashion.\nSome messengers permit being logged on from multiple (locations) computers at once and allow the current chat record to be fully available at all locations. This requires the instant messenger to support synchronous logins from multiple devices, however, which may be unwanted complexity to support, or may be considered a security issue. Without the extra complexity and security issues, it would be desirable to have a method to prevent messages from being untimely even if messages have been delivered to a computer that has been automatically logged off."}
{"text": "Over time or because of a particular event or condition (e.g., seismic activity, exposure to excessive or uneven loads or moments, exposure to micro-organisms, poor compaction, crown corrosion, corrosive soil, etc.), the structural integrity or capacity of force mains, other pipes and other structures may diminish. For example, such items may crack, corrode, deteriorate and the like. Different methods of repairing or otherwise strengthening damaged pipes and other items are well-known. For example, liners or layers made of fiber reinforced polymers can be attached to one or more portions of a pipe interior. Among other ways, liners of this type can be formed in a tubular shape and dragged in or everted into the pipe to be repaired. Still further, in larger diameter pipes, the lining may be formed by manually adhering sheets of repair material to the walls of the pipe.\nCertain high-strength liners or sheets are made of woven fibers of glass, carbon, aramid, or the like impregnated with a polymeric matrix that is, in some embodiments, cured in place. The woven fibrous material typically has a relatively open or loose weave that is water permeable in the absence of the polymeric material. Though the polymeric material in most instances prevents groundwater from infiltrating a pipe for a time, pinhole leaks are known to form in the openings between the woven fibrous strands under certain conditions. Pinhole leaks let undesirable groundwater enter the pipe. Moreover, high strength liners are often used to line pipes that carry pressurized fluids. Pressurized fluid can easily escape through pinhole leaks and cause further damage to the liner."}
{"text": "1. Field\nThe embodiments discussed herein are related to illuminating devices.\n2. Description of Related Art\nIn mobile devices, such as mobile phone handsets, ornamental elements are provided on a face panel or similar outer material. Light sources are used as one such ornamental element. In some cases, the face panel may be replaced to change the ornamental design. Such ornamental elements contribute to improved flexibility in product design and value (i.e., product appeal).\nTechnology related to such ornamental elements has been proposed, wherein a face panel is removably attached to a housing of a device. Light from light-emitting mechanisms on the housing is then made to illuminate the face panel, and a translucent display pattern on the face panel is displayed. (See, for example, Japanese Unexamined Patent Application Publication No. 2002-125016.)\nIn addition, technology has been proposed wherein an illumination case is removably attached to a main body of a mobile device. The illumination case is then made to light up by light from a light source on the housing, or by light from a through-hole formed on the main body. (See, for example, U.S. Pat. No. 6,435,690.)"}
{"text": "The present invention relates to a device for visually indicating that an object has been stored for an excessive period or above a predetermined temperature or for indicating that a package or seal has been opened or tampered with.\nNumerous devices have been proposed which provide a visual indication of the passage of a predetermined period of time, tampering, or storage under undesirable conditions. For example, Olsen et al, U.S. Pat. No. 3,344,670 (1967) describes a time/temperature indicator consisting of a cellulosic material such as filter paper treated with silver nitrate. This indicator exhibits a permanent color change depending solely upon the thermal history of the treated paper. Jackson, U.S. Pat. No. 3,480,402 (1969) discloses a time indicator formed of an absorbent carrier having absorbed thereon at least one chemical compound which changes color upon exposure to oxygen. The carrier and absorbed compound are protected from ambient oxygen by a barrier layer through which oxygen controllably diffuses over a preselected time. Thus, a color change is evident if recommended storage times are exceeded. Kydonieus et al, U.S. Pat. No. 4,212,153 (1980) discloses a time indicator in which migration of an agent from an interior layer to an exterior surface layer produces a visually perceptible color change. In one embodiment the migrating agent is a dye; in another it is an acid or base and the exterior layer contains a pH indicator. Other indicators are described in Yoshikawa et al, U.S. Pat. No. 4,169,811 (1979); Seiter, U.S. Pat. No. 3,967,579 (1976); and Youngren, U.S. Pat. No. 4,382,700 (1983)."}
{"text": "The present invention relates generally to material handling equipment and relates more particularly to pneumatic pop-up units adapted for use in material handling equipment.\nThe manufacture and/or packaging of many types of commercial articles often involves the use of material handling equipment. One example of material handling equipment is a bench or workstation at which certain assembly or packaging steps for an article are performed. Another example of material handling equipment is a conveyor system for transporting an article from one workstation to another.\nAs can readily be appreciated, an article that is seated on top of a bench or workstation often must be re-positioned or rotated thereon so that certain assembly or packaging steps can be performed. Similarly, it is often the case that the direction in which an article is traveling along a conveyor system must be altered or its movement altogether stopped. For these reasons, various omni-directional devices and stopping devices have been devised that are adapted for use in material handling equipment.\nFor example, in U.S. Pat. No. 4,627,526, inventor Masciarelli, which issued Dec. 9, 1986, and which is incorporated herein by reference, there is disclosed a conveyor system which includes a number of driven rollers mounted in a spaced, parallel relationship. A conveyor element is located between adjacent driven rollers, the conveyor element consisting of a main body having a series of passages in which plungers are slidably carried. An inflatable bladder moves each of the plungers from a retracted position to an operative position. More specifically, each passage is a bore in the main body, and each plunger consists of a sphere rotatably carried in a sheath whose outer surface is formed to slide smoothly in the bore. In this manner, when the bladder is inflated, the spheres are moved from their retracted positions to their operative positions, where they engage the article and lift the article above the top level of the driven rollers, thereby enabling the article to be re-positioned relative to the rollers (or pushed off of the driven rollers, for example, onto a lateral conveyor). When the bladder is then deflated, the spheres retract and an article situated thereon is lowered back onto the driven rollers. Alternatively, instead of including a sphere, the plunger may include a stop pin for use in stopping the movement of the article across the series of driven rollers.\nAs another example of an omni-directional device adapted for material handling equipment, in U.S. Pat. No. 4,660,994, inventor Masciarelli, which issued Apr. 28, 1987, and which is incorporated herein by reference, there is disclosed an anti-friction (i.e., omni-directional) assembly adapted for use in material handling equipment. More specifically, the aforementioned assembly comprises a main body having a first bore entering one surface and a second bore concentric with but non-coextensive with the first bore. A piston is slidable in the second bore and has a stem whose free end is slidable in the first bore. The stem carries a ball which is rotatable on a nest of secondary balls in a cavity in the free end of the stem. The piston and the stem are movable from a first position in which the ball lies substantially within the first bore below a friction element (e.g., the workstation, driven rollers, etc.) to a second position in which the ball protrudes from the first bore above the friction element. Means are provided for introducing fluid under pressure to the second bore to move the piston and to carry the ball from the first position to the second position. In addition, means are provided in the second bore to bias the piston to normally maintain the piston and the ball in the first position.\nAs still another such example., in U.S. Pat. No. 6,019,211, inventor Masciarelli, Jr., which issued Feb. 1, 2000, and which is incorporated herein by reference, there is disclosed a conveyor having a plurality of spaced rollers mounted between a pair of elongated horizontal frame members for supporting an article and enabling the article to be conveyed longitudinally of the frame members. The conveyor has a workstation formed by a plurality of interrupter strips extending in the spaces between rollers. Each interrupter strip has a plurality of lifting pads which are selectively movable between a lower position in which the upper surface of the pads are below the tops of the rollers and an upper position in which the upper surface of the pads are above the tops of the rollers for supporting the article and frictionally engaging the article to maintain the article in a stable position above the strips. The interrupter strips also have a plurality of balls mounted for omnidirectional rotation on the interrupter strips. Each ball is selectively movable between a lower position in which the top of the ball is below the top of the rollers and an upper position in which the top of the ball is above the top of the rollers. The balls are selectively moved independently of the lifting pads so that the article can be supported solely by the balls or by the lifting pads.\nReferring now to FIGS. 1 through 4, there are shown various views of still yet another example of a conventional pneumatic pop-up unit, said unit being of the omni-directional variety and being represented generally by reference numeral 11. Unit 11 comprises a housing 13, housing 13 being a generally cylindrical, cup-shaped member having a circular side wall 15, a bottom wall 17 and an open top. A nipple 19 is mounted in bottom wall 17, nipple 19 being adapted for connection to a pressurized gas (or other fluid) supply. Housing 13 is typically made of machined metal (e.g., aluminum, steel).\nUnit 11 also comprises a cap 21, cap 21 also being made of machined metal (e.g., aluminum, steel). Cap 21 comprises a top wall 23 and a side wall 25. Top wall 23 is flat and annular in shape. Side wall 25, which extends downwardly from the inside diameter of top wall 23, is circular in cross-section. Side wall 25 is dimensioned to fit closely inside side wall 15 of housing 13, with top wall 23 extending radially outwardly across side wall 15 and defining a circumferential flange 27 used to mount unit 11 in a hole H formed in a workbench W or the like. A pair of spaced apart circumferential lips 29 and 31 are formed on the outside surface of side wall 25 of cap 21, lips 29 and 31 defining a groove 33 therebetween. The outside surface of side wall 25 tapers slightly inwardly in diameter from lip 31 to the bottom surface thereof. Housing 13 and cap 21 are secured to one another by staking (i.e., mechanically deforming) housing 13 into groove 33 at approximately four points 34 equidistantly spaced around the circumference of housing 13.\nUnit 11 further comprises a piston 41, piston 41 being slidably mounted within housing 13. Piston 41, which is made by stamping and rolling a suitable sheet of metal, is a generally cylindrical, hollow member and includes a bottom wall 43, a circular side wall 45 and an open top. A ring 47 is fixedly secured by a press-fit to the outside surface of side wall 45, ring 47 being disposed a short distance from bottom wall 43. The top surface of ring 47 is used to engage the bottom end of a spring 49, the top end of spring 49 being engaged by lip 31 of cap 21. Spring 49 serves to bias piston 41 downwardly towards bottom wall 17 of housing 13. A rubber seal 51 is snugly fitted around side wall 45 of piston 41 and is engaged at its top surface by the bottom surface of ring 47 in such a way that ring 47 keeps seal 51 from sliding up on piston 41. Seal 51 serves to create an air-tight chamber between the bottom end of piston 41 and housing 13, seal 51 typically being of the split design variety (e.g., U-cup, block-V) so that, when seal 51 is subjected to upward fluid pressure, its outer leg is urged against the inside surface of wall 15 of housing 13 and its inner leg is urged against the outside surface of wall 45 of piston 41.\nUnit 11 further comprises a bowl-shaped race 55, race 55 being disposed inside of piston 41. Race 55 is seated on top of a race support 57, race support 57 being seated directly on top of the inside surface of wall 43 of piston 41. Race 55, which is made by stamping, has an outer diameter slightly smaller than the inner diameter of piston 41 and is adapted to hold a plurality of ball bearings 59. Race support 57 is shaped to include a circular side wall 61, on which race 55 is seated, and an annular bottom wall 63, which extends radially outwardly from side wall 61 to slightly less than the inner diameter of piston 41. A primary ball 65 is seated on ball bearings 59 and extends partially through the open top of piston 41, the top of side wall 45 being, rolled around ball 65 to retain ball 65 within piston 41. Ball 65 is aligned with the central opening in cap 21 and is appropriately sized relative thereto to partially extend therethrough.\nAlthough unit 11 has been found to be generally satisfactory for its intended purpose, the present inventor has observed the following shortcomings in connection therewith: First, as noted above, housing 13 and cap 21 are made by metal machining processes (either by screw machining or by powdered metal technology). Unfortunately, however, these metal machining processes often result in tolerances that are greater than optimally desired to enable housing 13 and cap 21 to fit together properly on a consistent basis when both are mass-produced. In addition, the aforementioned machining processes can be time-intensive and labor-intensive and, therefore, costly.\nSecond, as explained above, housing 13 and cap 21 are secured to one another by staking housing 13 into groove 33 of cap 21 at four equally spaced points 34 around the perimeter of housing 13. Unfortunately, however, said staking technique is, on occasion, not secure enough for cap 21 and housing 13 to withstand becoming separated by the force of the pressurized gas admitted into housing 13 to raise piston 41 therewithin and/or by the force of articles dropped or loaded on top of unit 11 during routine material handling steps.\nThird, once cap 21 has been staked to housing 13 in the above-described manner, cap 21 cannot be removed from housing 13 and thereafter re-attached without causing damage to cap 21 and/or to housing 13 (see, for example, breaks 36 in lip 31 of cap 21, as shown in FIG. 4). This is unfortunate since the internal components of a unit can become worn and/or dirty over time. Consequently, at present, the only viable option is to replace the entire unit.\nFourth, as described above, the upper end of spring 49 wraps around side wall 25 of cap 21 and is abutted at its top surface by lower lip 31. Consequently, when piston 41 is elevated to its uppermost position within housing 13, spring 49 becomes completely compressedxe2x80x94an undesirable condition since it may result in spring 49 weakening or locking in its compressed state. Moreover, because the outside surface of side wall 25 tapers inwardly from lip 31 to the bottom of side wall 25, the diameter of the wire used to make spring 49 must be small enough to allow spring 49 at its top end to fit between the rather small space left between side wall 25 of cap 21 and housing 13. Typically, this amounts to a wire diameter of about seventy-thousandths of an inch for spring 49.\nFifth, the manner in which piston 41 is finished, with race 55, race support 57, bearings 59 and ball 65 all arranged therewithin is as follows: First, piston 41 is stamped into the shape of a hollow cylindrical can having an open top. Race support 57, race 55, bearings 59 and ball 65 are then positioned within piston 41 in the manner described above. Finally, the top edges of side wall 45 of piston 41 are rolled or coined inwardly around ball 65 so as to entrap ball 65 and the other contents of piston 41 therewithin. Unfortunately, the metal used to form piston 41 has a tendency to xe2x80x9cspring-backxe2x80x9d to its original shape after having been stamped or rolled/coined; consequently, it is necessary to xe2x80x9coverbendxe2x80x9d the metal to compensate for anticipated spring-back. Such overbending, however, can create its own problems. For example, the overbending of the top edges of piston 41 around ball 65 can create nicks or cuts in ball 65 or can otherwise damage ball 65, particularly if ball 65 is made of a soft material, such as a soft plastic. In addition, the overbending of the top edges of piston 41 against ball 65 results in a considerable force being applied by race support 57 to the bottom of piston 41, thereby occasionally resulting in the bottom of piston 41 becoming deformed (i.e., spreading radially outwardly and/or buckling downwardly towards the middle). Furthermore, the overbending of the top edges of piston 41 against ball 65 can cause bearings 59 to be driven into and to deform race 55.\nIt is an object of the present invention to provide a novel pneumatic pop-up unit adapted for use in material handling equipment.\nIt is another object of the present invention to provide a pneumatic pop-unit as described above that overcomes at least some of the disadvantages associated with existing pneumatic pop-up units.\nIt is still another object of the present invention to provide a pneumatic pop-up unit that can be mass-produced readily and relatively inexpensively.\nAdditional objects, features, aspects and advantages of the present invention will be set forth, in part, in the description which follows and, in part, will be obvious from the description or may be learned by practice of the invention. In the description, reference is made to the accompanying drawings which form a part thereof and in which is shown by way of illustration specific embodiments for practicing the invention. These embodiments will be described in sufficient detail to enable those skilled in the art to practice the invention, and it is to be understood that other embodiments may be utilized and that structural changes may be made without departing from the scope of the invention. The following detailed description is, therefore, not to be taken in a limiting sense, and the scope of the present invention is best defined by the appended claims."}
{"text": "1. Field of the Invention\nThe present invention relates to a technology for determining an occurrence of a failure related to a head in a magnetic disk device that includes an automatic amplifying adjustment function.\n2. Description of the Related Art\nVarious techniques are know for preventing destruction or loss of data stored in the storage devices such as a magnetic disk device. For example, Japanese Patent Application Laid-open No. H9-139040 discloses a technology for preventing destruction or loss of data in the storage devices by assuring an operation to be properly performed for writing and reading of data to and from the storage devices.\nAccording to the conventional techniques, as shown in FIG. 13, when a magnetic disk device 1 (see, FIG. 11 and FIG. 12), which has been connected to a higher-level device via a host interface or the like, is performing a data writing processing, if a head 14 of the magnetic disk device 1 comes in contact with dust 25 on a surface 26 of a disk medium 15, the head 14 temporarily floats abnormally due to a shock from the contact and moves to a position described with a reference numeral 22. Accordingly, an extra distance shown with reference numeral 23 is generated between the head 14 and the surface 26 of the disk medium 15, and there causes a problem that data cannot be written onto the disk medium 15 because a magnetic field 24 of the head 14 cannot reach the surface 26 of the disk medium 15 so that a weak signal is output from the head 14. When a weak signal is output from the head 14, the processing is terminated to prevent the data destruction, and a new attempt is made to perform the data writing processing.\nAccording to the conventional techniques, floating of the head can be detected only after the occurrence of the floating. Thus, there is a need of a technology that can predict occurrence of the floating of the head."}
{"text": "A non-aqueous electrochemical cell using lithium as an active anode material has high energy density, good storage characteristics and wide operation temperature range. A non-aqueous electrochemical cell is therefore often used as a power source for a calculator, a watch or a memory back up system. Such a cell comprises an anode, an electrolyte and a cathode. In general, such a cell uses as an anode an alkali metal such as lithium or sodium; as an electrolyte or electrolytic solution, a solution of a solute such as lithium perchlorate or lithium tetrafluoroborate in a non-aqueous solvent such as propylene carbonate, .gamma.-butyrolactone or diglyme; and as a cathode, manganese dioxide or polycarbonmonofluoride.\nThe combination of relatively high theoretical energy density, potentially long life, and low cost materials such as reported in the sodium-sulfur system high temperature batteries is suitable primarily for low rate performance work such as electric road vehicle propulsion or load leveling of electric power supplies. The sodium-sulfur systems, first proposed in 1966, have had a great deal of effort expended in trying to develop a practical system. The basic operating principle involves the separation of two active molten materials, sodium and sulfur, by either a ceramic membrane of beta alumina or sodium glass, which at about 300.degree. C. or higher allows the passage of sodium ions that form with the sulfur any of the several polysulfides. The open circuit voltage of the system is at just over 2 volts, about the same as the lead-acid cell. Two formidable problems exist at the present time, viz., cracking of the separator and corrosion of the casing and seal.\nAnother somewhat similar system is the lithium-iron sulfide system, operating at about 450.degree. C. However, insufficient development has been done to date to demonstrate the widespread practicality of this system.\nAnother of the developments being pursued involves a lithium-based cell, in which the negative electrode is a lithium alloy (typically either lithium-aluminum or lithium-silicon), the positive electrode is an iron sulfide, and the electrolyte is a molten salt, such as the eutectic composition in the lithium chloride potassium chloride system. Because of the high melting point of such salts, such cells must be operated in the temperature range of 400-500 degrees centigrade.\nThis requirement to operate at such high temperatures has several important disadvantages. One of these is that various degradation processes, such as corrosion of the cell container, seals, and other components are accelerated by such high temperatures. Another is that a substantial amount of energy is lost through heat transfer to the surroundings. Still another is that the voltage obtained from such cells is lower at elevated temperatures, due to the fundamental property of the negative temperature dependence of the free energy of the cell reaction. Furthermore, the higher the temperature of operation, the greater the potential problems related to damage to the cell during cooling to ambient temperature and reheating, whether deliberate or inadvertent. Differences in thermal expansion, as well as dimensional changes accompanying phase changes, such as the freezing of the molten salt, can cause severe mechanical distortions, and therefore damage the cell components.\nCells involving a lower temperature molten salt electrolyte have been investigated where the molten salt is based upon a solution of aluminum chloride and an alkali metal chloride. However, the soluble positive electrode materials are not stable in the presence of the respective alkali metals. As a result, an auxiliary solid electrolyte must be used to separate the alkali metal and the counter electrode. One example of such a cell involves a molten sodium negative electrode, a solid electrolyte of sodium beta alumina, a molten aluminum chloride-sodium chloride salt, and either antimony chloride or an oxychloride dissolved in the chloride salt as the positive electrode reactant. Such a cell can operate in the temperature range of 150-250 degrees centigrade. It has the disadvantage of having to employ an electrolyte, which increases the cell impedance, as well as adding to the cost and complexity.\nU.S. Pat. No. 3,844,837 to Bennion et al discloses a non-aqueous battery in which the anode may be lithium and/or graphite on which lithium metal is deposited and as a positive electrode a platinum cup filled with powdered K.sub.2 SO.sub.4 and graphite is utilized. The electrolytes disclosed are LiClO.sub.4, LiCF.sub.3 SO.sub.3 or LiB.sub.4 dissolved in dimethyl sulfite.\nU.S. Pat. No. 4,877,695 discloses a non-aqueous electrochemical cell for use in primary rechargeable storage devices in which the cathode comprises an electrically conductive carbonaceous material, the anode is a molten mixture of two elements selected from the group consisting of sodium, potassium, cesium and rubidium, and the electrolyte comprises a solvent and an electrolyte salt selected from the group consisting of an alkali metal tetrafluoroborate and a tetraalkylammonium tetrafluoroborate. The present invention provides a specific improvement in the electrochemical cell of the patent by increasing its energy and power density.\nU.S. Pat. No. 4,886,715 discloses a primary rechargeable energy storage device comprising an anode which is an alkaline earth or alkali metal and a carbonaceous fiber cathode. The electrolyte is a membrane comprising lithium laurate."}
{"text": "1. Field of the Invention\nThis invention relates to beverage containers which have separate ice and beverage compartments. In such containers, the ice compartment is in heat exchange contact with the beverage compartment to keep the beverage cold while it is being consumed, without diluting the beverage as the ice melts.\n2. Description of Related Art\nMany beverages are best when served cold, including beer, soda and iced tea among others. Ice is the preferred way to keep the beverage cold as it is consumed, and for some beverages, such as soda and iced tea, it is common to place the ice directly into the beverage container where it dilutes the beverage as it melts. For other beverages, such as beer, this dilution is unacceptable.\nFor such beverages, even where the beverage and the beverage container are chilled prior to serving, the beverage rapidly becomes unacceptably warm even in an insulated container. The problem is compounded by the fact that many chilled beverages are consumed during the warmest months of year.\nTo eliminate the problem of dilution, various methods of cooling the beverage have been tried. One method has been to use a plastic-coated freezable gel material either as an artificial ice cube floating in the beverage, or as a strap-on coolant pack outside, but in contact with the beverage compartment. Both methods are expensive, and require additional freezer equipment and/or space for the numerous gel-cubes or gel-packs to be used. The artificial ice cube method has the added difficulties of risk of contamination from the gel material and sterilization costs in reusing the cubes.\nAnother method has been to use standard ice to cool the beverage, but to isolate the ice from the beverage in a separate compartment, thereby excluding the meltwater from the beverage compartment. This method has proven somewhat more satisfactory due to the existing investment and wide availability of ice-makers found in most homes and food service businesses.\nNonetheless, previous beverage containers employing the separate compartment design suffer from several deficiencies which have prevented their wide acceptance. One difficulty has been that most such containers fill both the ice compartment and the beverage compartment from the top. In this arrangement, the opening for the ice must compete for space with the opening for the beverage. The opening to the ice compartment is invariably too small making it difficult to fill with ice quickly. Enlarging the opening of the ice compartment makes it easier to fill with ice, but more difficult to fill with beverage.\nA related difficulty has been that it is too easy to inadvertently get ice in the beverage compartment or beverage in the ice compartment when filling the separate compartments. The difficulty of filling the ice compartment in prior designs is less objectionable in low volume applications, but in commercial establishments the time involved to carefully fill the ice compartment of each beverage container cannot be justified.\nEven those beverage containers which may be filled from the bottom, thereby eliminating the possibility of mixing beverage and ice, have not been satisfactory in commercial establishments. In some cases the beverage container is too fragile, too thin or is made of a material not suitable for scooping into the ice (such as glass, porcelain or other material which may chip). In other cases, the container has no handle, has a handle that is too thin and cannot be gripped tightly, or has a handle not positioned properly on the container. These problems prevent the container from being comfortably used in the inverted position to quickly scoop crushed ice or ice cubes into the ice compartment. Each of these deficiencies leads to a decrease in the speed with which the container can be filled.\nYet another problem has been that prior art designs do not cool the beverage efficiently. In normal use, the beverage container is initially filled with a pre-cooled beverage. Since a large percentage of the beverage will be immediately consumed, there is little time for the beverage to be warmed by the surrounding air and a relatively low ratio of cooling area to beverage volume is required. Later, with a partially full beverage container, the cooling ratio needs to be increased to keep the remaining beverage cool against the warming effect of the surrounding air. Prior art containers have not been designed to provide this variable cooling ratio.\nStill another difficulty with prior art designs is the relatively low heat transfer rate between the beverage and ice. When ice is directly mixed with the beverage, the large surface area of the ice results in rapid heat transfer. However, in prior art designs, the wall between the ice compartment and the beverage compartment has significantly less surface area than ice immersed in the beverage, resulting in a low heat transfer rate. Consequently, the beverage may not be cooled sufficiently rapidly to counteract heat transferred into the beverage from the surroundings.\nA final difficulty with prior art designs is that the ice and beverage compartments have usually been formed as separate components. This requires additional assembly and fastening or sealing operations during manufacture that greatly increase the cost of such containers and make them uncompetitive with simpler containers lacking the ice compartment.\nIn view of these deficiencies in the prior art, one object of the present invention is to provide an inexpensive beverage container with an ice compartment which can be easily handled when inverted and quickly filled by scooping in the crushed or cubed ice commonly found in most beverage serving commercial establishments.\nYet another object is to provide a beverage container with an ice compartment which can be sealed to inhibit melting of the ice. It is still another object of the present invention to provide an ice compartment which provides greater cooling relative to the remaining volume of beverage as the beverage is consumed.\nA further object of the invention is to provide a beverage container with an improved heat transfer rate between the beverage and the ice."}
{"text": "This subject matter concerns a dispensing system for a fluid in general, and more particularly a disposable inverted dispensing system for dispensing shampoo or other liquid materials (such as soap or conditioner, or gel or lotion of any type) for use in a shower.\nThe manufacture and sale of hair care and personal hygiene products constitutes an expansive market. As can be seen in almost any household, shampoo, conditioners, liquid soaps, and other materials represent a large portion of this market. Conventionally, such products are sold in plastic containers which stand upright with a capping device on top through which the fluid, for example, shampoo, conditioner, liquid soap, or other lotions is dispensed. The designs of such containers create several different drawbacks and problems from storage to dispensing. Such plastic containers litter the bathroom shower area of many households creating unsightly clutter. Often, the containers are not readily accessible within the shower area which can lead to the spilling of the container's content. Such clutter, inaccessibility and possible spills can create safety hazards within the shower area, as well as create waste and other disadvantages.\nThe upright containers also create a dispensing issue due to the configuration of the container in which it is contained and/or the potentially viscous nature of the content being dispensed. The nature of these containers requires the user to perform several steps to dispense the contents of the container. The user typically must pick up the container, invert it and usually squeeze it by applying the user's own pressure on multiple sides of the container to extract the fluid of the container. Further, with the dispensing opening of these containers on the upper end of the container when it is in a stored position, the fluid settles in the bottom of the container prior to use. The more the fluid in the container is used, the longer the distance the fluid must travel with less momentum due to the decrease in the mass of the fluid after each use, and the greater the effort and time required to dispense what fluid remains in the container. Often, the user must shake or otherwise manipulate the container to obtain sufficient fluid from the container.\nOther devices and containers have tried to address these problems of the upright containers with varying degrees of success, but all have their own distinct disadvantages. Hanger mechanisms can be attached to upright containers or can hold the upright containers to allow them to be hung upside down on the shower rod or the neck of the showerhead, thereby removing the containers from the floor of the shower and allowing the contents of the containers to settle at the dispensing zone of the container. However, these hanger mechanisms require the user to spend time attaching the hanger mechanism to the container. Some require a hanger to be screwed into the container, creating the opportunity for inadvertent discharge of the fluid. Others require the container to be inserted into a holding mechanism which is then tightened around the container. After the content is dispensed from the container, the holding mechanism is then loosened and the container has to be removed from the hanger mechanism requiring more time consumption.\nAlso, such hanger mechanisms impede the shower curtain when hung on the shower curtain rod or crowd the neck of the showerhead which has become an often-used location to otherwise hang personal care items. Further, the user still must grab the container on multiple sides and squeeze to dispense the contents.\nPermanent soap dispensers as seen in many public restrooms may tend to alleviate certain aspects of clutter, but the user has the burden of filling the dispenser on a periodic basis. For these permanent dispensers, containers of shampoo, conditioner, liquid soap, or other liquid material must be purchased in a store and then the contents of these containers must be transferred into the permanent dispensers. The transfer of such fluid is time consuming and, if care is not taken, creates an opportunity to spill (i.e., waste) the fluid, which in many instances may be rather expensive, particularly for “premium” products such as designer hair care formulas. Further, permanent soap dispensers are affixed to a shower wall in such a manner that, upon removal, lasting marks are left on the wall where the dispenser had been attached.\nA need still exists for a dispensing system for shampoos, conditioners, liquid soaps or other consumable personal care fluids that does not add to the clutter around a shower area, is easily accessible in the shower, and decreases the possibility of spillage. A need also exists for the dispensing system to be easily assembled and disassembled to allow for the quick set up (i.e. installation) and disposal of the dispensing system, while still permitting the user to dispense the shampoo, conditioner, soap, or other liquid in a well facilitated fashion without having to shake or manipulate any part of the dispensing system."}
{"text": "In recent years, many medical appliances such as an endoscope, a medical manipulator, and the like, which have a tubular portion which is long and thin and a tip end thereof is inserted into a body cavity for observation or an operation, have been developed and used. Since the medical appliances are used by being inserted into a body cavity of a patient, from a hygiene viewpoint for prevention of infection, it is necessary to wash and clean both inside and outside of the medical appliances for them to be clean, by washing treatment such as washing, sterilizing, and rinsing, or with combination thereof, when re-using the medical appliances. However, since the medical appliances have long and thin tubes as those are necessary for functions and usages, a cleaning agent or a disinfectant does not sufficiently pass through the inside of the thin tubes during the washing treatment.\nA washing treatment apparatus and a washing treatment method of the related art for such medical appliances having a thin tubular structure are disclosed in Patent Literature 1, for example. The washing treatment apparatus of the related art includes a washing bath in which an endoscope is loaded, a connection tube to which at least a washing solution is supplied, a guide member which guides the connection tube to an axial direction of a channel port of the endoscope provided on an external side of the washing bath, a movement mechanism which moves the connection tube through the guide member, a sealing member for sealing the connection tube and the channel port, and a partition which forms a sealed structure in which the movement mechanism and a part of the connection tube are provided."}
{"text": "1. Technical Field\nThe present invention relates to semiconductor gain media, especially to those devices that include an integrated multimode region and single mode region and more particularly to laser devices with single spatial mode, diffraction-limited emission with flared gain regions.\n2. Background Art\nExternal-cavity semiconductor lasers, including lasers with frequency selective tuning elements in the cavity, are well known and have been extensively studied. For example, T. Fujita, et al., in Applied Physics Letters 51(6), pages 392-394 (1987), describe a laser having a buried heterostructure laser that has been antireflection (AR) coated on the intracavity facet, a collimating lens, a polarization beamsplitter, external cavity mirrors in each of the TE and TM polarization light paths, and an electro-optic modulator in the TE polarization path between the beamsplitter and cavity mirror. The configuration allows selection of either the TE or TM mode of oscillation by adjusting the modulator's bias voltage. W. Sorin, et al., in Optics Letters 13(9), pages 731-733 (1988), describe a laser having a laser diode with one of its facets AR coated to reduce its reflectivity, a lens, a single mode optical fiber and a tunable evanescent grating reflector for providing feedback. The laser is wavelength tunable by sliding the feedback grating laterally over the fiber. P. Zorabedian et al., in Optics Letters 13(10), pages 826-828 (1988), describe another wavelength tunable laser using either a rotatable interference filter in an external Fabry-Perot cavity or an external grating reflector providing frequency-selective feedback.\nA problem with previously available external-cavity semiconductor lasers is their generally low output power (on the order of 10 mW cw and 200-300 mW pulsed). Further, higher output powers are associated with unstable output intensity and frequency and less than good modal quality.\nIn U.S. Pat. No. 4,251,780, Scifres et al. describe semiconductor injection lasers that are provided with a stripe offset geometry in order to enhance and stabilize operation in the lowest order or fundamental transverse mode. In one configuration, the stripe geometry has a horn shaped or trapezoidal section connected to a straight section, in which the width of the horn shaped or trapezoidal section expands from 8 .mu.m at the straight section to 25 .mu.m at the cleaved end facet. In contrast to configurations in which the edges of the stripe waveguides are linear and orthogonal to the cleaved end facets of the lasers, the nonorthogonal angled or curved edges of the offset stripe geometries cause higher order modes to reflect or radiate out of the waveguide, thereby increasing the threshold of the higher order modes relative to the fundamental mode.\nIn U.S. Pat. No. 4,815,084, Scifres et al. describe semiconductor lasers and laser arrays in which lenses and other optical elements have been integrated into the semiconductor bodies of the lasers by means of refractive index changes at boundaries in the light guiding region, where the boundaries are characterized by a lateral geometric contour corresponding to surfaces of selected optical elements so as to cause changes in shape of phase fronts of lightwaves propagating across the boundaries in a manner analogous to the change produced by the optical elements. In one embodiment, a biconcave or plano-concave diverging lens element is integrated within the laser in order to counteract the self-focusing that usually occurs in broad area lasers and that can lead to optical filamentation and lateral incoherence across the laser. The diverging lens in the laser allows the laser to operate as an unstable resonator, leading to high output power and good coherence across the lateral wavefront.\nAn object of the invention is to provide a high power, semiconductor gain medium which emits a single spatial mode, diffraction-limited output beam.\nAnother object of the invention is to provide a semiconductor gain medium chip that includes a multimode region and a single spatial mode region with an optical cavity that includes a cleaved facet reflector and a grating reflector."}
{"text": "Conventionally, as a known technology, at the time of switching from self-driving to manual driving, an attentiveness of a driver is detected, and the switching operation is performed according to the detected attentiveness of the driver."}
{"text": "1. Field of the Invention\nThe present invention relates to the structure of a compact ceramic multi-layer substrate for dealing with high frequency signals.\n2. Description of the Prior Art\nPreviously, there have been known monolithic substrates integrating active and passive elements on a semiconductor substrate and hybrid substrates integrating various active and passive elements on a dielectric layer such as alumina. In general, the hybrid substrate has a single layer structure. A multi-layer structure of electric circuits is also known for a dielectric substrate made of alumina or the like.\nIn the above-mentioned structure of a monolithic substrate, single layer structures are integrated. Therefore, if the sizes of electronics components to be integrated are determined, the total size can be determined depending on how the components are packed in the same plane. However, for a high frequency circuit, electromagnetic coupling exists through the space or the substrate. Therefore, it is difficult to integrate a circuit at a density more than a certain value.\nFurther, for a hybrid substrate, the structure may be made from multi-layers. However, only metallization is made from multi-layers, so that electronics components have been put on the surface. Therefore, this limits the integration and the compactness.\nFurther, in general, a high frequency circuit is liable to be coupled electromagnetically through a space. In order to prevent such coupling, various techniques such as a metallic shield plate arranged at the top are needed to shield the electromagnetic coupling. Then, high frequency characteristics may be damaged if the profile becomes complicated, if the size becomes large or if the shield is incomplete.\nIn order to solve these problems, various structures of ceramic multi-layer substrates have been proposed. For example, T. Nishikawa et al. proposed a compact high frequency filter circuit composed of passive components such as inductances (L) and capacitances (C) by using a multi-layer structure made from dielectric ceramics and electrodes (RF Front End Circuit Components Miniaturized Using Dielectric Resonators for Cellular Telephones, IEICE Transactions, Vol. E 74, No. 6, (1991), pp. 1556-1562). This structure is made from the same dielectric material. However, for a circuit such as a resonator including active elements for example transistors, it is not appropriate to use a first dielectric material of large dielectric constant on a substrate section wherein discrete components are set because a phase shift occurs at high frequencies. On the other hand, it is preferable that the dielectric constant of a second dielectric material for constructing the resonator is larger than that of the former dielectric material.\nIn order to solve this problem, for example, K. Hayashi discloses a structure wherein the dielectric constant of a first dielectric material for constructing a resonator is larger than that of a second dielectric material in a section for setting discrete components (Japanese Patent laid open Publication No. 79,601/1992). However, in a circuit including active components such as transistors, a capacitor of large capacitance such as a bypass capacitor is needed for a direct current bias circuit. Therefore, this structure is not sufficient for miniaturization. Further, in a circuit such as a high frequency amplifier, a path for direct current bias is connected to a high frequency transmission line. Then, if high frequency components pass the direct current bias circuit, unfavorable phenomena such as a loss or unstableness inducing oscillation may occur.\nIn a structure disclosed in Japanese Patent laid open Publication No. 25,082/1991, the diameter of a throughhole of a multi-layer substrate connected to an electrical conduction line for power supply is decreased or the length thereof is increased to increase an inductance component, so as to avoid high frequency components in the connection line for power supply. However, if the diameter is decreased, the amount of the current is limited, while if the length is increased, the miniaturization is difficult."}
{"text": "It is known that electrical noise occurs in large scale integration (LSI) semiconductor circuits because of, for example, the presence of alpha (α) rays or neutron beams that cause the semiconductor circuits to malfunction. The a rays are emitted when radioisotopes included in the packages, wiring, etc. of the LSI semiconductor circuits have decayed, and the neutron beams originate from cosmic rays. The above malfunction is called a soft error in contrast to a hard error caused by a failure of hardware in a semiconductor circuit or the like. Accordingly, the soft error differs from the hard error in that the soft error causes transient malfunction in which the operation of the semiconductor circuit can be recovered while the hard error causes permanent destruction at a specific part in the circuit.\nHowever, if the transient malfunction caused by the soft error is stored in a latch circuit making up the LSI semiconductor circuit, the malfunction caused by the soft error may cause a fixed error in the semiconductor circuit.\nThe increasing integration and decrease in size of LSI semiconductor circuits reduce the charge capacities of storage nodes of latch circuits included in the LSI semiconductor circuits to increase the probability of soft errors caused by occurrences of electrical noises becoming fixed errors. This is because the logical states of the storage nodes of the latch circuits are apt to reverse due to the electrical noise to increase the frequency of the malfunction states.\nUnder such situations, approaches to improve the soft-error resistance of logic LSI circuits including latch circuits have been proposed in recent years. For example, a method of determining the output by a decision of the majority of signals output from three or more latch circuits receiving the same data is proposed. However, with this method of avoiding errors, the number and area of circuits are increased because of the three or more latch circuits, a majority circuit, and so on to cause a reduction in the circuit performance. In addition, such a latch circuit in the related art can keep the normal output if a soft error occurs at a single node but cannot keep the normal output if soft errors continuously occur before new data is input. This is because the output is determined on the basis of the majority rule.\nIn order to reduce the effects of the above problems, a latch circuit including multiple circuit units and feedback paths through which and the multiple circuit units the input is fed back to the input is known. Complementary metal oxide semiconductor (CMOS) inverter circuits are generally used as the circuit units. In each of the CMOS inverter circuits, inverter circuits, particularly, P-channel transistor (PTr) and N-channel transistor (NTr) are coupled in series to each other.\nAccordingly, it is desired to realize a latch circuit having an improved soft-error resistance."}
{"text": "This patent relates to the field of devices for indicating loss of pressure and imminent failure of pneumatic tires in vehicles.\nU.S. Pat. No. 2,686,493 shows, mounted to the fixed axle portion of the vehicle, at the axle-leaf spring connection, a sounding wire, consisting of a curved wire extending vertically downward to make a resonant sound through the frame of the vehicle upon contact to the wire with the ground.\nThis is extended by U.S. Pat. No. 2,794,411, which improves this concept by providing a resonator, a metal cup and an adjustment mechanism for varying and fixing the length of the contact wire.\nAnother ground contact detector is a sound generating device which is wheel mounted. Reeder U.S. Pat. No. 1,601,604 to Reeder shows an attachment in the hubcap of the vehicle in which a single striker mechanism is depended radially outward along the wheel of the vehicle. Contact with the ground causes it to strike and ring an internal bell mechanism.\nU.S. Pat. No. 1,785,614 is typical of the majority of patents for tire pressure loss detectors in that it utilizes a mechanism physically mounted upon the wheel which responds to the outward expansion of a tire carcass during deflation. All such patents, of course, represent mechanisms whose behavior cannot be predicted in the event of a sudden failure, such as a blowout, with the attended collapse of the tire structure. Boyer shows a spring supported mechanism which is triggered from an inner to an outer position, if the tire wall extends to an outward point. The outward motion is not visually detectable but is rather used to trigger a microswitch creating an electrical signal for an electric signaling circuit.\nU.S. Pat. No. 902,693 is an earlier mechanism for responding to the outward expansion of the tire. A spring loaded arm swings to an outward position to create a sound signal by striking against a vibrating member located on the frame of the vehicle.\nU.S. Pat. No. 1,947,760 is a frame rather than wheel mounted deflection indicator of this basic type; the outward deflection of the tire pulls a cable mechanism which then triggers a mechanical sensor. This device is mounted upon the fixed portion of a brake drum of a car, and it would appear to be subject to high failure rate in the event of a strike by road debris and the like.\nU.S. Pat. No. 2,917,020 shows an alternate form: a hubcap mounted sound alarm in which spring wires within a hubcap are set so that they will strike the ground during deflation, creating a ringing effect on the hubcap. An apparently unpatentable version is shown in FIGS. 8 through 9, which appear to be a duplicate of an earlier patent to Neiswanger.\nU.S. Pat. No. 3,937,542 shows an electric low pressure sensor with indicator in which an electric pressure sensor is mounted within the valve stem of a tire and coupled to a motor vehicle indicator mounted through an insulated slip spring and brush mechanism on the frame of a car.\nPerry U.S. Pat. No. 3,742,896 discloses a pneumatically actuated signal flag mounted upon the axle hub of a vehicle and responsive to the pneumatic pressure to indicate a low pressure condition in a vehicle tire."}
{"text": "The delivery of text-based messages (i.e. data packets) from a sending device to one or more receiving devices over a wireless LAN, presents special challenges. Typically, the message is routed through a wireless gateway where it is stored until it has been transmitted to, and stored within, an electronic mail server of the data network. Receiving devices are then able to retrieve stored messages from the electronic mail server at their convenience. The speed at which electronic messages are transmitted from a sending to a receiving device depends in part on how efficiently data packets are transported from a sending mobile device to an electronic mail server through wireless communication networks.\nWhen a wireless gateway receives a data packet from a mobile device over a wireless network, the received data packet is sent to a destination electronic mail server. However, to ensure that the data packet is successfully transmitted to the destination electronic mail server and not lost in the meantime, the wireless gateway generally stores the data packet in an internal permanent storage device (e.g. a database server or a file system) before transmitting the data packet. Typically, the wireless gateway waits until the permanent storage device confirms storage of the data packet before proceeding with processing the data packet or even with sending acknowledgement of the receipt of the data packet back to the mobile device. This kind of storage procedure appreciably slows down the processing of data packets within the router."}
{"text": "This disclosure relates generally to tracking devices, and more specifically, to locating a tracking device or when the tracking device is in a communicatively-restricted environment, such as out of range of mobile devices.\nElectronic tracking devices have created numerous ways for people to track the locations of objects. For example, a user can attach a tracking device to an important object, such as keys or a wallet, and use the features of the tracking device to more quickly locate the object (e.g., if it becomes lost). In particular, the tracking device can emit a tracking signal that can be detected by one of a community of mobile devices. The location of a community mobile device that is within range of the tracking device and receives the tracking signal can be used to estimate the location of the tracking device. This location can be reported to the user's mobile device so that the user can find the lost object.\nAs community mobile devices move around, they enter and exit the ranges of tracking devices. However, if none of the community of mobile devices come within range of a lost tracking device, the user may be unable to find a lost object. Thus, current tracking devices are limited by the range of the tracking devices and the coverage provided by the community of mobile devices.\nThe figures depict various embodiments of the present invention for purposes of illustration only. One skilled in the art will readily recognize from the following discussion that alternative embodiments of the structures and methods illustrated herein may be employed without departing from the principles of the invention described herein."}
{"text": "The G protein coupled nociceptin receptor known as ORL1 has been shown to be involved in the modulation of pain in animal models. It bears high homology to the classic opioid receptors (μ, k, δ), but has little cross reactivity with their native ligands. Current opioid analgesics target these classic opioid receptors, but have limiting side effect profiles (e.g. tolerance, physical dependence, respiratory depression and decrease of gastrointestinal function). ORL1 receptors are colocalized in regions of the CNS similar to the opiod receptors, as well as in the periphery.\nNociceptin, the endogenous ligand to ORL1, was discovered in 1995 and shown to be a peptide ligand that activates the ORL1 receptor, but not the classic opioid receptors. Initial reports have suggested that nociceptin and the ORL1 receptor are involved in a newly discovered pathway involved in the perception of pain. Further reports have shown nociceptin to be analgesic when administered intrathecally to rodents. The in vivo efficacy of nociceptin in animal models of pain is similar to that of the endogenous opioids. Nociceptin is also reported to act as an anxiolytic agent when administered into the brains of rodents. The in vivo efficacy in rodent anxiety models is similar to classic benzodiazepine anxiolytics. In addition, nociceptin has been recently reported to inhibit capsaicin induced bronchoconstriction in isolated guinea pig lung tissue, suggesting a role for ORL1 agonists in the treatment of cough. Together, these data suggest that nociceptin receptor agonists may have significant analgesic, anxiolytic, or antitussive properties."}
{"text": "Level shifters, which receive and convert input signal(s) of smaller signal range to output signal(s) of larger signal range, are important blocks of interface circuits. For example, in gate driver chip for driving display panel, internal control signals inside the chip have a signal range of 0 to 3 Volts, but signals outputted to gates of display panel require a signal range of −8 to 3 Volts. For conversion between two different signal ranges, a level shifter is adopted to convert an input signal range of 0 to 3 Volts to an output signal range of −8 to 3 Volts.\nPlease refer to FIG. 1 illustrating a conventional shifter 10. The conventional level shifter 10 includes a pair of (p-channel MOS, Metal-Oxide-Silicon) transistors TP1, TP2 and a pair of (n-channel MOS) transistors TN1 and TN2. An input signal IN is inverted to another input signal INB by an inverter INV, wherein the input signals IN and INV have a signal range between voltages VPP and VSS. The level shifter 10 operates between voltages VPP and VGL for respectively providing output signals OUT and OUTB at nodes n2 and n1 according to the input signals IN and INB, with a signal range of the output signals OUT/OUTB expanded between voltages VPP and VGL. Gates of the transistors TP1 and TP2 respectively receive the input signals IN and INB, and gates of the transistor TN1 and TN2 are respectively coupled to the nodes n2 and n1.\nIn the prior art, the operation of level shifter 10 can be described as follows. For example, when the input signal IN transits from the voltage VPP to the voltage VSS, the transistor TP1 starts to turn on, and then pulls the voltage of the node n1 (i.e., the output signal OUTB) up to the voltage VPP. As the voltage of the node n1 raises, the transistor TN2 starts to turn on, and then pulls the voltage of the node n2 (the output signal OUT) down to the voltage VGL.\nHowever, when the input signal is originally maintained at the voltage VPP, the transistor TN1 is turned on. When the input signal IN transits from the voltage VPP to the voltage VSS, the turned-on transistor TP1 has to compete against the turned-on transistor TN1. The turned-on transistor TN1 tends to keep the node n1 to the lower voltage VGL, so the transistor TP1 has to conduct more current than the transistor TN1 with better conduction for pulling the voltage of the node n1 up to the voltage VPP.\nSince the transistor TP1 is a p-channel transistor, it suffers from lower current driving ability (e.g., lower carrier mobility). Therefore, a larger aspect ratio (ratio of channel width to length, W/L) is required to drive enough current against the transistor TN1. This requirement leads to a disadvantage that the layout area of the transistors TP1 and TP2, as well as the total layout area of the level shifter 10, can not be reduced.\nIn addition, when the transistor TP1 competes against the transistor TN1 owing to transition of the input signal IN, the large current conducted by the transistor TP1 causes a larger short current of longer duration which continues to drain power from the voltage VPP during the competition, and increases power consumption of the conventional level shifter 10.\nFurthermore, in the prior art, the level shifter 10 is driven by voltage. That is, the transistors TP1 to TP2 and TN1 to TN2 operate respectively according to cross voltages between their gates and sources. Thus, circuit design must be specifically customized according to voltage ranges of the input and output signals. The circuit design has to be tailored for different applications. For example, if an application requires a level shifting from an input signal range of 0 to 3 Volts to an output signal range of −3 to 3 Volts, and another application needs a level shifting from an input signal range of 0 to 3 Volts to an output signal range of −20 to 3 Volts, the level shifter designed for former application can not be adopted for the latter application in the prior. The level shifter has to be re-designed with different parameters (e.g., larger circuit dimensions) to fit the latter application."}
{"text": "The current global framework, with growing oil price instability, limited oil and gas resources and the Kyoto Protocol environmental requirements, calls for continued improvements in the usage of renewable energy resources, including solar. During the last decades, interest in solar energy solutions has increased because the potential of solar energy has become apparent.\nEven though solar radiation is a source of high temperature and energy at the origin, sun-earth geometrical constraints lead to a dramatic dilution of flux and to irradiance available for terrestrial of about 1 kW/m2 and consequently, to a supply of low temperature to the thermal fluid. It is therefore essential requisite for solar thermal power plants to make use of optical concentration devices that enable the thermal conversion of high solar flux and with relatively little heat loss. Typical concentration devices for solar radiation are: parabolic troughs, linear Fresnel reflectors, parabolic dishes, and power tower (or central receiver solar systems—CRS). Solar thermal power based on such concentrators is called Concentrated Solar thermal Power (CSP). The most common solar thermal power systems are parabolic trough and power tower.\nA parabolic trough power plant is basically composed of three main elements: the solar system, the steam generator and the power block. The solar system comprises parabolic trough collectors' field and the oil (i.e. heat transfer fluid) circuit. At the solar field, solar radiation is collected and converted into thermal energy as the temperature of the oil circulating through the receivers, increases. Once heated at the solar field, the oil is pumped and conveyed to a steam generator, to produce high pressure/temperature steam. The steam is then used to produce power similarly to conventional steam power plants (e.g. Rankine cycle).\nIn power towers (or CRS), incident sunrays are reflected from large mirror collectors (heliostats), which track the apparent sun movement and concentrate the energy flux onto a solar receiver, where energy is transferred to thermal fluid. This fluid can also be direct steam. The thermal energy conversion into electricity is quite similar to fossil-fueled power plants.\nBoth parabolic trough and CRS are proven technology. However, their cost is not yet competitive with conventional alternatives unless subsidized. Therefore, it is yet desired to improve the existing plants cost effectiveness.\nU.S. Pat. No. 5,444,972 describes a hybrid power plant driven by hydrocarbon fuels and solar energy. The solar system is CRS and the power block is condensing steam cycle (Rankine). Exhaust gas of a turbine and fuel added burner backup the solar field in generating the steam. Sensible heat storage and a receiver thermal shield are described in this publication as being optional features.\nGB 2449181 discloses a combined gas/steam (or Brayton/Rankine) hybrid power plant. The heat input to the gas cycle partially depends on solar heating of air supplied to an initial combustion stage from a compressor of the gas turbine engine. Simplification and/or greater responsiveness of the gas turbine engine's fuel control system is achieved by providing the power plant with separate combustion stages being coupled to an associated turbine. Compensation for variations in insolation is achieved by controlling fuel flow to the initial combustion stage while compensation for variations in electrical load is achieved by controlling fuel flow to a subsequent combustion stage. The final turbine stage drives a generator, and the exhaust from the final turbine stage is passed through a heat recovery steam generator to raise further power in a steam turbine that also drives a generator. The second combustion stage, including a superheater, may be located in the steam cycle part of the combined cycle to further heat the steam before entry to the steam turbine. Concentrated solar energy may heat a heat transfer fluid, e.g. a molten salt, for intermediate heat transfer and heat storage.\nU.S. Pat. No. 5,417,052 describes a hybrid, combined cycle power plant driven by CRS and natural gas. The combined cycle is Brayton/Rankine. The CRS heat transfer fluid (HTF) is molten salt. The molten salt preheats the air leaving the gas cycle compressor. The exhaust gases from the gas cycle drive a steam cycle for additional energy production.\nUS 2008127647 discloses a method for oversizing the Rankine power block with respect to the topping cycle block in combined cycles where the bottoming cycle is Rankine. The method describes the production of the additional steam in the Rankine plant by solar energy. The method is applicable in plants that are equipped with an oversized heat recovery steam generator (HRSG) and steam turbine system that is coordinated with duct burners or other means for additional steam generation.\nUS 20060260314 discloses an integration between a combined cycle power plant and a solar Rankine power plant. Relatively high temperature, low pressure reheat from the combined cycle power generation system can be used, through, for example, a superheater, to raise the temperature and pressure of a working fluid in a solar Rankine power generation system. The resulting integrated system has enhanced efficiencies as compared with stand-alone systems.\nMost of the solar thermal power today is produced by parabolic trough technology. Typically, parabolic trough plants consist of large fields of parabolic trough collectors, a heat transfer fluid/steam generation system, a Rankine steam turbine/generator cycle, and optional thermal storage and/or fossil-fired backup system(s). The solar field is modular in nature and comprises many parallel rows of single-axis-tracking parabolic trough solar collectors, normally aligned on north-south horizontal axis. Each solar collector has a parabolic shaped reflector that focuses the sun's direct beam radiation on a linear receiver located at the parabola focal axis. The collectors track the sun from east to west during the day. A heat transfer fluid (hereinafter “HTF”) is heated up to about 400° C. while being circulated through the receiver and returns to a series of heat exchangers (hereinafter “HE”) in the power block where the solar heat absorbed by the HTF (typically synthetic oil) is extracted to generate high pressure superheated steam (e.g. 100 bar, 371° C.). The superheated steam is then fed to a conventional reheat steam turbine/generator (condensing turbine) to produce electricity. The expanded steam from the turbine is condensed in a saturated condenser and returned to the HE via feed-water pumps to be transformed back into steam. Wet cooling towers or sea water supply cold water to the condenser. After passing through the power block, the cooled HTF is re-circulated through the solar field.\nThe parabolic trough approach is currently the most proven and lowest cost large-scale solar power technology available today as described by H. Price et al. in “Advances in Parabolic Trough Solar Power Technology”, Journal of Solar Energy Engineering V. 124, Issue 2, pp. 109-125, May 2002. However, the cost of parabolic trough power is not yet competitive with conventional alternatives unless subsidized (typical installation cost for 50 MW plants is in the range of 4-6 $/W). One of the reasons is that the standard condensing steam power block described above has several drawbacks that highly limit the overall plant cost effectiveness and deployment. Some of these drawbacks are the following:\n1) Water Cooled Condenser:\nThe condensers in steam power units are typically water cooled type where the water are typically cooled by a wet cooling tower (WCT). Air cooled condensers (dry cooling towers—DCT) are inferior because they increase the steam condensation temperature, and the steam condensing cycle efficiency is very sensitive to that temperature. In addition, DCT are more expensive than WCT. Consequently, solar thermal power plants are installed mostly in locations where water is available, rather than in desert and arid areas where solar radiation levels are typically the highest.\n2) Condensing Turbine:\nCondensing steam turbines for solar plants, typically designed to generate 50-80 MW, have a very unique design which creates low market availability (up to two years of manufacturing lead time). This might introduce a significant bottleneck to projects time tables because the other components lead times are substantially lower (typically less than one year). Another downside of the condensing turbine is the performance penalty when scaling down. For instance, the difference between 60 MW and 5 MW turbine shaft (isentropic) efficiencies is about 6%. The performance penalty increase both the specific required collection area (m2/W) and the electricity specific cost ($/W). This is critical because the small and medium size plants (5-20 MW) market is currently growing rapidly.\n3) Operation Time (Capacity Factor):\nDuring day time, the solar radiation intensity varies. In a clear sky day, the radiation level increases gradually during the morning and reaches maximal level around noon while from afternoon to sunset, the insolation declines. In a standard (condensing turbine) power block, steam is fed into the turbine only when reaching a certain desired temperature, typically around 370° C. and an initial flow rate. Operating the turbine with substantially lower steam temperatures is eluded because that would expose the turbine blades to high levels of moisture which is destructive. Therefore, the standard routine in such a solar plant would be to wait until the heat transfer fluid heats up to about 390° C. so the desired steam top temperature can be achieved. However, this oil temperature can be obtained only when reaching a minimum insolation level, say 400 W/m2. Thus, during early mornings and late afternoons' periods the solar plant cannot produce power. The result is shorter plant operation time (or lower capacity factor) and of course, higher electricity cost.\n4) Working Fluid:\nSteam condensation increases the power block complexity and cost due to the following two reasons:\na) The saturated steam specific volume is very high (for example, 12 m3/kg at 50° C.), requiring large size condensing unit.\nb) The saturated pressure of steam at typical condensing temperatures is much below atmospheric (for instance, at 35° C. the pressure is 0.056 bar). Consequently, the condenser must always be maintained under deep vacuum and free of oxygen. Addressing these constrains, increases the plant O&M costs.\nOther power cycles are also described in the literature.\nFor example, A. I. Kalina in “Combined-cycle system with novel bottoming cycle”, Journal of Engineering for Gas Turbines and Power v. 106, pp. 737-741 (1984) describe a thermodynamic cycle that is designed to replace the Rankine cycle as a bottoming cycle for a combined-cycle energy system as well as for generating electricity by using low-temperature heat sources.\nP. A. Lobos and E. D. Rogdakis described in their publication “A Kalina power cycle driven by renewable energy sources”, Energy, Jan. 31, 2009, a Kalina cycle using low-temperature heat sources to produce power. The main heat source of the cycle is provided by flat solar collectors. In addition, an external heat source is connected to the cycle, providing 5% up to 10% of the total thermal energy used in the cycle. The cycle operates at low pressure levels (0.2-4.5 bar) and low maximum temperature (130° C.). For given conditions, an optimum range of vapor mass fractions and operating pressures can be identified that result in optimum cycle performance. Simple equations have been derived that link the operational parameters with the independent variables as well as with the cycle efficiency.\nThe present invention seeks to provide a solution that eliminates the major part of the above problems while increasing the solar plant efficiency."}
{"text": "In order to restore or assume the function of a damaged or lacking tissue or organ the transplantation of natural tissues or organs from another donor or, if possible, the concerned individual himself is an established practice that has long been known. Due to the constant lack of suitable donor tissues and donor organs and other disadvantages of natural tissue, e.g., rejection reactions and the risk of the transfer of diseases from the donor to the recipient, many efforts are being directed to the production of artificial tissues and organs as alternatives.\nA customary procedure for this is the provision of a carrier or matrix structure from material compatible with the body that is colonized with differentiated cells of the target tissue, and the cultivation of the cells in vitro until a tissue-like cell structure has been produced. The differentiated cells are obtained either from cultures of explanted tissue samples or from stem cells that had been stimulated to differentiate. The use of stem cells permits a more rapid production of larger amounts of the desired cells in many instances. Traditionally, pure cell populations of a certain type are produced. Most of the in vitro organs or in vitro tissues known in the state of the art are disadvantageous in as far as that they do not have or do not develop the tissue structure that corresponds to the morphological constitution of the native tissue or organ even after implantation and a fairly long residence time in the body. This applies as a rule even when the carrier matrix had been colonized with several different populations of tissue-typical cells.\nAnother approach for treating degenerative diseases or damage to tissues and organs using stem cells consists in implanting the stem cells and/or differentiated cells derived from them directly into the damaged tissue/organ in order to proliferate there and result in a regeneration of the damaged tissue/organ. In this instance too the implantation or regeneration of a certain cell type traditionally is in the foreground. The problem is, however, that in very many degeneration phenomena or damages to various organs a plurality of cells is always involved so that, e.g., in the skin keratinocytes, epithelial cells and blood vessels are also affected in addition to fibroblasts. In the case of nerve damage even glia cells often have to be replaced too and muscle damage often means destruction of the accompanying nerves."}
{"text": "This application is based upon and claims the benefits of priority from the prior Japanese Patent Application No. 2001-100402 filed on Mar. 30, 2001; the entire contents of which are incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to an infrared sensor and a manufacturing method thereof and, in particular, to a pixel construction of an uncooled infrared sensor and a manufacturing method thereof.\n2. Related Art\nInfrared image sensing is characterized in its ability to pick up images at night as well as during the daytime and its higher transmittance to smoke and fog than visible radiation. Being capable of acquiring temperature information of an object to be sensed, infrared image sensors are adaptable to a variety of applications, such as a monitor camera and a fire detection camera, besides the field of defense.\nIn recent years, extensive studies have been made on xe2x80x9cuncooled infrared solid-state imaging elementsxe2x80x9d which eliminate the need for the cooling mechanism to allow for a low-temperature operation, which is the most critical shortcoming of a quantum type infrared solid-state imaging element or a leading conventional type. In such an uncooled infrared solid-state imaging device, an incident infrared ray having a wavelength of 10 xcexcm or so is converted into heat by an absorber and then the temperature change caused by this small amount of heat at a heat sensing portion is converted into an electric signal by thermoelectric conversion means of any kind, before reading out the electric signal to obtain infrared image information.\nThere are three ways of improving the sensitivity of such an uncooled infrared sensor, generally classified as follows:\nThe first method is to improve the ratio of an infrared power incident to an infrared detection portion, dP, to a temperature change of an object, dTs, namely dP/dTs. This method achieves improvement mainly by optical systems, which corresponds to enlargement of an infrared lens aperture, application of an antireflective coating, use of a lens material of low optical absorption, increase of an infrared absorptivity at an infrared detection portion or increase of an infrared absorption area. In accordance with the recent multiplication of pixels for an uncooled infrared sensor, the size of a unit pixel has mainly come to as small as 40 xcexcmxc3x9740 xcexcm and, among the items mentioned above, the increase of an infrared absorption area at an infrared detection portion has remained a relatively critical problem to be solved. However, a report has been published that an infrared absorption area can be increased up to 90% of a pixel area by stacking an infrared absorption layer on top of the pixel (Tomohiro Ishikawa, et al., Proc. SPIE Vol. 3698, p. 556, 1999) and it is difficult to obtain a further considerable amount of improvement in sensitivity by any optical means.\nThe second method is to improve the ratio of a temperature change at an infrared detection portion, dTd, to an incident infrared power, dP, namely dTd/dP. Whereas the aforementioned method is an optical one, this can be said to be thermal. Generally speaking, in an uncooled infrared sensor to be packed in a vacuum package, transportation of heat from an infrared detection portion to a supporting substrate is predominantly accomplished by heat conduction by a supporting structure which supports the infrared detection portion in a cavity within the supporting substrate. Thus, attempts have been made to lay a leg-like supporting structure made of a material of low thermal conductivity as narrowly and longways as possible within the constraint of the design. See Tomohiro Ishikawa, et al., Proc. SPIE Vol. 3698, p. 556, 1999, for example. As shown in FIGS. 14(a) and 14B, a construction has been made such that a pair of narrow, spiral slits 201 is formed around a pixel 200 to form a cavity 204 at the bottom of the element. These spiral channels 202 are used as a supporting structure to support the pixel 200 afloat and wirings 203 are provided to connect to a peripheral circuit. When a pixel is being miniaturized to as small as 40 xcexcmxc3x9740 xcexcm or so, however, since fine processing has already been made at the silicon LSI processing level, a considerable degree of further improvement in sensitivity may hardly be realized through refinement of the layout of the supporting structure. Similarly, it is difficult to further reduce the thermal conductivity which is a material characteristic of the supporting structure. In particular, with regard to the wiring for outputting electric signals from the infrared detection portion, it is difficult to realize a considerable amount of improvement in sensitivity in terms of material since there is a requirement contradictory between the electric conduction and heat conduction whose mechanisms are similar.\nThe third method is to improve the ratio of an electric signal change caused by thermoelectric conversion means, dS, to a temperature change at an infrared detection portion, dTd, namely dS/dTd and is, therefore, electrical in nature. This method is, unlike the other two, is directed to sheer sensitization, that is, improvement of dS/dTd; however, it is quite important to reduce various electric noises generated. Thus, various kinds of thermoelectric conversion means have so far been investigated. Reported are, for example, a thermopile for converting a temperature change into a potential change by means of Seebeck effect (Toshio Kanno, et al., Proc. SPIE Vol. 2269, pp. 450-459, 1994), a bolometer for converting a temperature change into a resistance change by means of the temperature change of a resistor (A. Wood, Proc. IEDM, pp. 175-177, 1993), a pyroelectric element for converting a temperature change into an electric charge by pyroelectric effect (Charles Hanson, et al., Proc. SPIE Vol. 2020, pp. 330-339, 1993), a silicon pn junction for converting a temperature change into a voltage change by certain forward currents (Tomohiro Ishikawa, et al., Proc. SPIE Vol. 3698, p. 556, 1999), etc.\nIf a comparison is made among those methods, however, no one method is then decisively superior to the others in comprehensive view of their thermoelectric conversion characteristics, noise characteristics and manufacturing methods. For instance, the bolometer is superior in terms of temperature resolution while the silicon pn junction, which can be manufactured by the silicon LSI processing alone, is superior in terms of manufacturing method.\nAs described above, one of the methods for sensitizing an uncooled infrared sensor is a secondly mentioned thermal one, which improves the ratio of a temperature change at an infrared detection portion, dTd, to an incident infrared power, dP, namely dTd/dP. Generally speaking, transportation of heat from an infrared detection portion to a supporting substrate is predominantly accomplished by heat conduction through a supporting structure which supports the infrared detection portion in a cavity within the supporting substrate. Thus, attempts have been made to lay a leg-like supporting structure made of a material of low thermal conductivity as narrowly and longways as possible within the constraint of the design. When a pixel is being miniaturized to as small as 40 xcexcmxc3x9740 xcexcm or so, however, since fine processing has already been made at the silicon LSI processing level, a considerable degree of further improvement in sensitivity may hardly be realized through refinement of the layout of the supporting structure.\nFurthermore, with a trend of miniaturizing pixels and supporting structures involved in the development of fine processing technology in silicon LSI processing, the influence of the heat transportation by emission from the bottom of a pixel and supporting structure will predictably be appreciable and the sensitization only through the reduction of heat conduction by the miniaturization of a supporting structure will predictably be restricted by a limit in sensitivity contributable to such emission.\nIn addition, in order to read signals out from an infrared detection portion, a supporting beam line to be formed within a supporting structure must be provided with a unique wiring layer made of a low thermal conductivity material. For example, a supporting beam line structure using titanium material is known to be effective.\nExcept for the case where titanium itself is used as a thermoelectric conversion material, however, in addition to the wiring process at an element peripheral circuit, another process will be required for forming a supporting beam line structure alone.\nThe present invention is directed to improve the above supporting structure and to provide a highly sensitized infrared sensor and a manufacturing method thereof.\nAccording to one aspect of the present invention, provided is an infrared sensor device comprising:\na semiconductor substrate having a principal plane including a surface portion in which a hole is provided;\nan infrared detection pixel disposed in the hole, comprising an infrared absorption portion configured to absorb an infrared ray and convert it into heat and a thermoelectric conversion portion configured to convert the heat into an electric signal;\na semiconductor peripheral circuit, disposed on the principal plane of the substrate, configured to read out the electric signal from the infrared detection pixel; and\na supporting beam line with U-shaped cross section, extending across the infrared detection pixel and the surface portion having the hole, configured to support the infrared detection pixel, afloat within the hole and to electrically connect the infrared detection pixel to the semiconductor peripheral circuit.\nAccording to another aspect of the present invention, also provided is an infrared sensor device comprising: a semiconductor substrate having a principal plane including a plurality of surface portions in each of which a hole is provided;\na plurality of infrared detection pixels each disposed in the hole, each comprising an infrared absorption portion, configured to absorb an infrared ray and to convert it into heat and a thermoelectric conversion portion configured to convert the heat into an electric signal;\na semiconductor peripheral circuit, disposed on the principal plane of the substrate, configured to read out the electric signals from the infrared detection pixels, the semiconductor peripheral circuit including at least a MOS transistor having a gate comprising a first electric conductor with U-shaped cross section along the directions of the source and drain and a metal filled within the electric conductor; and\na plurality of supporting beam lines, each extending across the infrared detection pixel and the surface portion having the hole, configured to support the infrared detection pixel, afloat within the hole and configured to electrically connect the infrared detection pixels to the semiconductor peripheral circuit, each of the supporting beam lines including at least a second electric conductor with U-shaped cross section and being formed in the same layer as the first electric conductor with U-shaped cross section of the gate of the MOS transistor.\nAccording to still another aspect of the present invention, further provided is an infrared sensor device comprising:\na semiconductor substrate having a semiconductor supporting substrate and a semiconductor layer and an insulating layer interposed therebetween, in which a plurality of holes are provided on the semiconductor substrate, the holes penetrating the semiconductor layer and the insulating layer to reach the semiconductor supporting substrate;\na plurality of infrared detection pixels each disposed in each of the holes, each comprising an infrared absorption portion configured to absorb an infrared ray and to convert it into heat and a thermoelectric conversion portion configured to convert the heat into an electric signal;\na semiconductor peripheral circuit, formed on the semiconductor layer, including at least a MOS transistor having a damascene metal gate formed for driving the infrared detection pixels to read out the electric signal; and\ndamascene metal supporting beam lines, extending across the infrared detection pixels and side surfaces of the holes of the semiconductor substrate respectively, configured to support the infrared detection pixel, afloat within the holes respectively and to electrically connect the infrared detection pixels to the semiconductor peripheral circuit respectively, the damascene metal supporting beam lines being disposed in the same layer as the damascene metal gate.\nThe U-shaped electric conductor may have a shape of a square having one side thereof open, a circle or oval having a segment thereof open or a parabola.\nBy forming the supporting beam line across the infrared detection pixel and the substrate on the same layer as the damascene metal gate of the MOS transistor of the peripheral circuit, the beam line with U-shaped cross section may be rigid and elaborate, thereby considerably reducing the cross section of the supporting beam line.\nIt is therefore possible to considerably reduce the heat conduction through the supporting beam line, which predominates the heat transportation between the infrared detection pixel and the substrate, consequently providing a highly sensitive uncooled infrared sensor device.\nAccording to another aspect of the present invention, provided is a method of manufacturing an infrared sensor, comprising:\nembedding a silicon oxide layer in a substrate comprising a single crystalline silicon supporting substrate, an silicon oxide film layer formed on the single crystalline silicon supporting substrate and a single crystalline silicon layer formed on the silicon oxide film layer, over a predetermined region of the single crystalline silicon layer;\nforming a thermoelectric converter pn junction for an infrared detection pixel in the single crystalline silicon layer;\nforming a supporting beam line including a U-shaped electric conductor on the silicon oxide layer while forming a gate electrode of a MOS transistor having an electric conductor with U-shaped cross section, the MOS transistor being included in a peripheral circuit on the single crystalline silicon layer;\nforming an infrared absorption layer on the single crystalline silicon layer; and\nforming a hole by etching for isolating the infrared detection pixel of the substrate from the substrate and suspending the infrared detection pixel within the hole by the supporting beam line.\nBy forming the supporting beam line with a minimum cross section through the step of etching the hole for isolating the periphery of the infrared detection pixel from the substrate, the length of the supporting beam line may considerably be lessened and the aperture ratio, defined as an infrared detection pixel area relative to a unit pixel area, may considerably be improved. In addition, lessening the length of the supporting beam line may increase the mechanical strength and stabilize the operation while enabling sensitization to an infrared ray. Furthermore, linearization of the supporting beam line may extremely stabilize the manufacturing method thereof and, in turn, improvement of the yield may lower the cost.\nAt the same time, according to the present invention, the chip area may be reduced by miniaturizing the peripheral circuit and, therefore, the cost may naturally be lowered.\nAs described above, according to the present invention, an inexpensive and highly sensitive uncooled infrared sensor may be obtained."}
{"text": "Geologic formations may contain reservoirs of sought-after hydrocarbons. Boreholes are typically drilled into the reservoirs where the hydrocarbons can be extracted through rock fractures in the formations. Logging tools may be conveyed through the boreholes in order to characterize the formations so that drilling and/or completion resources can be used efficiently. Hence, innovations that improve the accuracy of formation characterization would be well received in the drilling and production industries."}
{"text": "A superacid is a known class of acidic material that has an acid strength, measured by Hammett acidity function H.sub.0, greater than that of 100% H.sub.2 SO.sub.4, which has an H.sub.0 value of -12. A superacid, therefore, has an H.sub.0 value of less than -12 or an acid strength greater than -12. Superacids are useful for reactions that are generally catalyzed by strong acid, such as paraffin isomerization.\nM. Hino and K. Arata describe the synthesis of a solid superacid having an acid strength of up to H.sub.0 .ltoreq.-16.04 by exposing hydroxides or oxides of Fe, Ti, Zr and Hf, prior to crystallization, to sulfate ions, followed by calcination in air at over 500.degree. C. in J. Chem. Soc., Chem. Commun., 1148 (1979). K. Arata and M. Hino also describe the synthesis of a solid superacid having an acid strength of H.sub.0 .ltoreq.-14.52 by impregnating Zr(OH).sub.4 or amorphous ZrO.sub.2 with aqueous ammonium metatungstate, followed by calcining in air at 800.degree. to 850.degree. C. in J. Chem. Soc., Chem. Commun., 1259 (1988) and in \"Synthesis of Solid Superacid of Tungsten Oxide Supported on Zirconia and Its Catalytic Action\", Proceedings 9th International Congress on Catalysis, Volume 4, pages 1727-1734 (1988).\nThe superacid described by K. Arata and M. Hino in \"Synthesis of Solid Superacid of Tungsten Oxide Supported on Zirconia and Its Catalytic Action\" is particularly useful as a catalyst in the isomerization of butane to isobutane and of pentane to isopentane. However, in order to obtain maximum catalytic activity, K. Arata and M. Hino report calcination temperatures of from 800.degree. to 850.degree. C. for the tungsten-modified zirconia catalyst.\nSoled et al describe, in U.S. Pat. No. 5,113,034, a solid acid catalyst having an H.sub.0 value ranging from -14.5 to about -16.5, comprising a sulfate or tungstate-modified Group IVB oxide, which is useful to dimerize C.sub.3 or C.sub.4 containing feedstreams. A calcination temperature range for the tungstate-modified zirconia catalyst of 450.degree. to 800.degree. C. is given. Specifically, the Examples use an initial calcination temperature of 600.degree. C., followed by a calcination temperature of 800.degree. C. prior to charging.\nU.S. Pat. No. 5,198,403 (Brand et al, herinafter Brand) discloses a catalyst for the selective reduction of nitrous oxide with ammonia which contains, in addition to titanium oxide as component A, at least one oxide of tungsten, silicon, boron, aluminum, phosphorus, zerconium, barium, yttrium, lanthanum, cerium as component B1 and at least one oxide of vanadium, niobium, molybdenum, iron and copper as B2, with an atomic ratio between the elements of components A and B from 1:0.001 up to 1. The catalyst of Brand requires the presence of titanium oxides, presumably because it is effective in the reduction of nitrous oxides. Reduction of nitrous oxides is the intended purpose of Brand. The catalysts of the instant invention do not contain titanium and are not drawn to nitrous oxide reduction.\nU.S. Pat. No. 4,918,041 (Hollstein et al, hereinafter Hollstein) is directed to a sulfated calcined solid catalyst. The catalysts of the instant invention are not sulfated. Group VI B metals such as tungsten are employed in the superacid, rather than Group VI A metals such as sulfur. Group VI B metals are preferred in the instant invention due to their stability and ease in regeneration.\nThe generation of acid activity in solid oxide catalysts in general, and in the tungsten-modified zirconia catalyst specifically, requires calcination of the catalyst at temperatures of about 800.degree. C. This extreme temperature, however, causes significant loss of catalyst surface area. For example, K. Arata and M. Hino report surface areas of 35.3 and 29.5 m.sup.2 /g for catalysts calcined at 800.degree. and 900.degree. C. in \"Synthesis of Solid Superacid of Tungsten Oxide Supported on Zirconia and Its Catalytic Action\". By contrast, the surface areas for catalysts calcined at 600.degree. and 700.degree. C. are reported as 44.2 and 38.5 m.sup.2 /g. Hence, it is clear that as the calcination temperature is increased, the surface area of the catalyst is decreased. Further, the extreme calcination temperature required to generate acid activity results in a more difficult manufacturing process."}
{"text": "Materials, such as wall patching compound or wall board joint compound are available, generally, in large containers such as tubs or drums, and smaller containers, such as tubes. Particularly, wall patching compounds are available in the smaller, tube-type containers.\nIn use, these materials are squeezed or expelled from the tube in an appropriate amount for the required task. The material is then spread over the area to be covered using a separate trowel or spreading knife.\nTypically, this operation requires the use of both hands to open the container, dispense or expel the material from the container, close the container, and spread the material on the surface which needs to be repaired. This operation can be particularly cumbersome when the user is standing, for example, on a ladder or an elevated surface.\nAnother drawback to known containers is that many such containers do not provide an air-tight seal to keep the material fresh and usable. The material which is exposed to air tends to harden and become unusable. This is particularly troublesome with tube-type containers which have small nozzle openings to discharge the material. If the material hardens in the nozzle, the tube is generally unusable, and the entire tube must be discarded.\nThus, there continues to be a need for a container for materials such as wall patching compound, which permits \"one-handed\" operation, and which container provides an air-tight seal to maintain the compound in a fresh, usable condition."}
{"text": "I. Field\nThe present disclosure relates generally to extension of battery life in wireless communication systems, and more particularly to extension of battery lifetime in Bluetooth equipped wireless communication devices.\nII. Background\nBattery lifetime is a chief concern in portable Bluetooth 802.11 computing devices, such as Ultra Mobile Personal Computers (UMPCs) and Personal Digital Assistants (PDAs). It is known that battery lifetime can be extended by keeping the host processor asleep whenever possible. It is useful for such computer devices to be aware when they enter 802.11 coverage, and wake-up the host processor, if needed, thereby eliminating the need for the user to manually check for coverage.\nIn this connection, wake-up conditions must be sophisticated, varying from the existence of a network to complex conditions such as the existence of an urgent e-mail, availability of a particular kind of service, etc.\nHowever, traditional wake-up solutions require the existence of alternative interfaces other than 802.11 to initiate wake-up, the presence of full-blown additional processors, or host processor hierarchies for wake-up initiation. Further, some solutions restricted to 802.11-based wake-ups are extremely inflexible (such as WoWLAN and AMD Magic packet). Moreover, traditional host processor wake-up conditions are not variable from application to application, which often change according to complex user requirements, such as the need to be notified of urgent e-mails, existence of a particular service, existence of a particular network and so forth. Current 802.11-initiated wake-up methods are very rigid and limiting.\nU.S. Pat. No. 7,398,408 disclose methods for accessing computing devices over wireless local area networks and more particularly to systems and methods for waking computing devices from a powered down or sleep state with signals sent over wireless local area networks. This is accomplished by:\nbroadcasting a signal containing a wake-up data sequence for at least one computing device over a wireless network, wherein broadcasting the signal comprises periodically broadcasting the signal over one or more wireless channels until either a first predetermined time period expires or a confirmation signal is received;                receiving the broadcast signal at the at least one computing device, while that device is in a reduced power mode; wherein receiving comprises first entering a wake-up data sequence detection mode;        scanning the received signal for a wake-up data sequence for that computing device, wherein scanning is performed until either the wake-up data sequence for that computer device is received or a second predetermined time period expires; and        restoring that computing device to a full power mode upon detection of a wake-up data sequence for that device and exiting the wake-up data sequence detection mode if the second predetermined time period expires.        \nThere is a need for an infrastructure system where 802.11 ATs are only awake when interaction with the AP is needed and possible, and asleep (in standby) otherwise.\nThere is also a need to provide a flexible mechanism whereby the AT's 802.11 interface can detect when “interaction with the AP Is needed” as possible."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of displaying technology, and in particular to a method for manufacturing a curved liquid crystal panel.\n2. The Related Arts\nA liquid crystal display (LCD) is composed of an enclosure, a liquid crystal panel arranged in the enclosure, and a backlight module mounted in the enclosure, in which the liquid crystal panel is composed of a color filter (CF) substrate, a thin-film transistor (TFT) array substrate, and a liquid crystal layer filled between the two substrates. The CF substrate and the array substrate have inside surfaces that face each other and comprise transparent electrodes formed thereon. The liquid crystal display applies an electric field to control the orientation of the liquid crystal molecules so as to change the polarization state of light and uses a polarizer to select transmission or blocking of a light path so as to achieve the purpose of displaying. With the progress of the displaying technology, people's demands of consumer electronics are no longer limited to functionality and also require satisfactory aspects in regard to design, aesthetics, and excellent visual experience. A curved liquid crystal display that has been created for such purposes is not only a novel design but also providing better visual effect and experience.\nA curved liquid crystal panel is a key component of a curved liquid crystal display. A conventional way of manufacturing a curved liquid crystal panel is to first laminate prepared CF substrate and array substrate in a flat condition to form a flat panel like liquid crystal panel and to subject the flat panel like liquid crystal panel to the application of an external force so as to make the liquid crystal panel that is originally in the form of a flat panel curved, where a degree of curving (radius of curvature) is determined according to a design size and viewing distance of the liquid crystal panel.\nIn a liquid crystal panel of a main stream LCD, a black matrix (BM) is provided on the CF substrate side to prevent the occurrence of light leakage at the site of each of pixels. Further, photo spacers (PSs) are provided between the CF substrate and the array substrate to maintain the distance between the two substrates in order to prevent the occurrence of abnormality.\nAs shown in FIG. 1, in the conventional way of manufacturing a curved liquid crystal panel, when the liquid crystal panel is being curved, since the CF substrate 100 and the array substrate 200 are moved relative to each other, a relative displacement of the black matrix 110 with respect to TFT pixels occurs in a curving portion, resulting in change of aperture ratio and causing dark blocks occurring in eye vision during a displaying process. When the liquid crystal panel is further curved, the relative shift between the CF substrate and the array substrate 200 is further enlarged so that two adjacent color resist blocks may move to locations close to a data line 210. The electric field around the data line 210 would cause the liquid crystal molecules to rotate and eventually, light leakage results. Such a displaying defect is even more apparent in displaying a pure color image. Thus, the relative shift between the CF substrate 100 and the array substrate 200 would cause movement of color resist, leading to color shift of the mixed color and irregularity of displayed color.\nAs shown in FIG. 2, due to the curving of the liquid crystal panel, the photo spacers 400 that were originally distributed in a reasonable manner in a flat panel are caused to generate different pushing effects that are generally in proportion to the different degrees of curving from the center of the liquid crystal panel to the two opposite, left and right, edges. With the curving of the liquid crystal panel, the curvature at the locations marked B, C, D and B′, C′, D′ in FIG. 2 changes so that photo spacers 400 at different locations would be of different heights. This readily leads to a potential risk that the cell thickness of the liquid crystal layer 300 becomes inconsistent due to the photo spacers 400 at different locations having different heights.\nIn brief, for a liquid crystal panel of which a circumference is adhesively fixed by enclosure resin 600, to change from a flat panel form into a curved condition would cause the array substrate 200 and the CF substrate 100 to move relative to each other and thus changes of the aperture ratio and the cell thickness of the liquid crystal layers 300 at different locations of the liquid crystal panel, in the entirety thereof and consequently displaying defects associated therewith would result."}
{"text": "The use of Digital Subscriber Line (DSL) technology depends heavily on the length and quality of the copper wire between the modem of an end user and the DSL Access Multiplexer (DSLAM). Longer copper wiring reduces the line quality and thus the bandwidth usable on this line. Reducing the length of the copper wire by locating a DSLAM near the end user, typically within a few kilometers, and serving only a small amount of users, for instance 120, will improve the service to the end user significantly.\nA DSLAM that is located near the end user is protected from weather conditions by some sort of enclosure, for instance an outdoor cabinet. The DSLAM requires a high bandwidth connection to the Central Office (CO) which could be an office of a telecom operator or internet service provider, and has multiple low bandwidth connections to end users generally used for xDSL access. Other technologies such as Integrated Services Digital Network (ISDN) or Plain Old Telephone Service (POTS) work similarly where a group of end users is connected to an Access Multiplexer (AM), which in turn provides access to the network of an operator.\nA Remote Unit (RU) is an enclosure which may be provided with at least one Access Multiplexer, one or more high bandwidth media to a CO and one or more low bandwidth links to end users. Such a Remote Unit is usually located above ground in what is referred to as street cabinet or outdoor enclosure. Alternatively a remote unit can be located underground in what is typically referred to as a manhole.\nA remote unit typically has several connections to the central office. Optical fiber provides a reliable high-speed medium which is ideal for carrying multiple information flows simultaneously. Putting fibers in the ground is however an expensive task with major disturbance for the public and traffic. Roads and sidewalks need to be opened before fiber can be placed underneath. Therefore spare fibers are made available in the remote unit for future expansion of the remote unit.\nA brochure concerning the Alcatel 7300 ASAM (Advanced Services Access Manager) ETSI version, Release 4.3/4.4 describes the use and advantages of the Alcatel 7300 ASAM in xDSL technology. Particularly the section titled “Extending Network Coverage” on page 6 of the brochure describes how the Alcatel 7300 ASAM can be used as a remote unit to extend geographical coverage of an xDSL deployment. It further describes how such a 7300 ASAM can be connected to products such as a micro DSLAM or a remote Alcatel 7300 ASAM-o as outdoor cabinet. Both the 7300 ASAM and subsequent connected devices such as a micro DSLAM can be considered a remote unit for the current invention. This brochure can be retrieved from the Internet via:\nhttp://www.alcatel.be/doctypes/opgproductbrochure/pdfa4/7300_etsi_bro.pdf.\nThe datasheet containing specifications of the Alcatel 7300 ASAM-c compact (ETSI version) release 4.3 discloses a DSLAM that can be included in an enclosure to form a remote unit. The Alcatel 7324 Remote Unit is yet another product offering DSLAM technology close to the end user premises. This datasheet can be retrieved from the Internet via:\nhttp://www.alcatel.be/doctypes/opgdatasheet/pdfa4/7300c_etsi_R43_ds.pdf.\nAll of the above-described products can be included in outdoor cabinets in order to be located near a particular amount of users. These outdoor cabinets provide protection from weather conditions and prevent unwanted physical access to the equipment inside.\nSeveral manufacturers provide enclosures for communication equipment. For instance ICEE Ltd. (Installers of Communication and Electrical Equipment) offers the design and installation of custom-made enclosures. Knürr AG provides a series of outdoor enclosures for various purposes in their Tecoras product line. The CS product line of Rittal is another example of outdoor enclosures that can be used to house communications equipment and be a part in the remote units described in this invention.\nContent broadcasters, such as television stations have a need to provide coverage of particular events, for instance music festivals or major news events. This coverage involves acquiring video and audio content, editing this content to an item ready for broadcast and transmitting this item to a central location for broadcast, not necessarily in that order. Typical parts in this process are satellite trucks which provide a satellite uplink to a central location and which are used to transmit the audio and video content. These trucks may also provide an editing studio to prepare the items before they are transmitted. Satellite uplinks are not always reliable; they are prone to delays and interference of radiation in the atmosphere. A typical example of this is the loss of a link during a live interview or the delays that occur between asking a question and a person actually hearing that question.\nIt is an object of the present invention to provide flexible means for delivering content to a central location from a remote location using existing network infrastructure. It is another object of the present invention to provide encoding and decoding means for easy content delivery. It is a further object of the present invention to facilitate real time audio and video communication between two remote points. It is another objective of the present invention to provide these means for delivering content in a secure way, restricting access from unauthorized persons or devices."}
{"text": "1. Field of the Invention\nThe present invention relates to a drive circuit and a drive method of a semiconductor laser module that externally modulates carrier light output from a semiconductor laser using an electro-absorption type optical modulator to output the modulated light, in particular, to a drive control technique for compensating for temperature characteristics of the semiconductor laser and the electro-absorption type optical modulator.\n2. Description of the Related Art\nIn various optical communication systems for transmitting optical signals over a long distance, there is sometimes used a semiconductor laser module in which a semiconductor laser and an external modulator of electro-absorption type are integrated with each other, since such a semiconductor laser module enables to increase a transmission distance compared with a semiconductor laser of direct modulation type. In the aforementioned semiconductor laser module, operation characteristics of the semiconductor laser and the electro-absorption type optical modulator (hereinafter referred to as EA modulator) each has dependence on temperature. Therefore, in order to hold optical output power in stable, it is necessary to stabilize a module temperature or to control a drive condition according to a temperature change.\nIn a conventional drive control of a semiconductor laser module having an EA modulator, such a method is typical in that, for example, by using an electronic cooling device such as a Peltier device and a temperature detection device such as a thermistor integrated with each other in a module, a current to be supplied to the Peltier device is controlled so that a resistance value of the thermistor is maintained to be constant, to hold the temperature of a semiconductor laser and the EA modulator to be constant. In this case, the EA modulator is driven by a constant modulated electric signal and at the same time, a constant drive current or a drive current to make a back power monitor current constant is supplied to the semiconductor laser, thereby achieving a stable optical output. The drive control method described above, however, has a disadvantage that since the Peltier device needs to be driven, the power consumption is increased and also a package size of the module is made to be large.\nTo cope with this problem, a semiconductor laser module using an EA modulator which is not provided with a Peltier device, has recently been under development. According to this semiconductor laser module, both the reduction of power consumption and the miniaturization can be achieved by eliminating the need of a Peltier device, although there is a need of stabilizing optical output power by controlling a bias voltage of the EA modulator or a drive current of the semiconductor laser according to a temperature change. Also, since this module is of an external modulation system using the EA modulator, it is possible to obtain an optical output having a small amount of chirping.\nAs a conventional drive control technique applicable to such a semiconductor laser module without the Peltier device as described above, a drive circuit for an EA modulator is disclosed in Japanese Unexamined Patent Publication No. 11-119176. In this drive circuit, an anode voltage of the EA modulator is detected and according to the detection result, a bias voltage of the EA modulator is controlled, so that an applied voltage to the EA modulator is maintained to be constant even if a temperature change or a change with age occurs.\nHowever, in the case where the semiconductor laser module provided with EA modulator without the Peltier device is driven by applying the conventional control technique as described above, although the bias voltage of the EA modulator is controlled according to the temperature change and the like, since the EA modulator is driven by a constant modulated electric signal, an optical output waveform of the semiconductor laser module is considerably deteriorated.\nFIG. 16 is a diagram showing temperature dependence of operating characteristics of a typical EA modulator. In FIG. 16, each characteristic curve represents a relation between power Pf of an optical signal output from the EA modulator and an applied voltage Vea to the EA modulator at each of temperatures 0° C., 25° C. and 75° C. In FIG. 16, if a modulated electric signal of a waveform as shown in the lower left part is applied to drive the EA modulator, since the optical output power Pf is changed along each characteristic curve, an optical waveform at 0° C. output from the EA modulator is considerably deteriorated as compared with an optical waveform at 75° C. as shown in the upper right part. Even if the bias voltage is controlled according to the temperature change as in the conventional technique so that the EA modulator operates in a region where the characteristic curve has the large inclination and is changed substantially linearly, since the inclination or the distortion of the characteristic curve in each operating region differs from each other depending on the temperature, an extinction ratio or a duty of the optical output waveform is changed depending on temperature when the EA modulator is driven by the modulated electric signal in which a modulated amplitude and a cross point are set to be constant.\nFurther, since a current generated by absorption of carrier light (hereinafter referred to as a photocurrent) flows through the EA modulator, there is caused a problem in that the applied voltage to the EA modulator deviates, if the photocurrent is changed when controlling a current source for drive controlling the EA modulator. Therefore, for the semiconductor laser module without the Peltier device, it is critically important to realize a drive control method that is not affected by a change in the photocurrent due to the electro-absorption effect, when the temperature change or the change with age is compensated by controlling the electrical drive signal of the EA modulator."}
{"text": "This patent document relates in general to on-chip waveguides with flattened dispersion, and in particular to on-chip generation of an ultra-broad optical spectrum, called a supercontinuum, based on the flattened dispersion of the waveguides.\nGenerating a supercontinuum from a narrowband pulse is of particular interest in both fundamental science and engineering research of modern optics. Although recent years have witnessed dramatic advances in supercontinuum generation using photonic crystal fibers, it is quite challenging to obtain an octave-spanning supercontinuum on a chip, partially because of strong chromatic dispersion in high-index-contrast nonlinear integrated waveguides."}
{"text": "Three-dimensional (3D) virtual worlds (for example, World Of Warcraft, Second Life, Runescape, etc.) allow users to enter and move about, as well as interact with other users and objects, within a three-dimensional virtual environment. Such 3D virtual worlds systems generally operate in the following manner:\n1) They require a 3D rendering engine to be resident on the client's system in order to render a user's viewport into the 3D world.\n2) The 3D rendering software further requires that the client machine have certain graphics processing capabilities which allow the 3D world to be rendered on the local machine. The more complex the world, the more powerful the client device's graphics capabilities must be.\n3) In order for a user to view a scene on her local system, the 3D rendering engine must have access to the assets necessary to build the scene accurately. Such assets typically include geometry mesh data (vertices, index and attribute buffers), material and shader definitions (including procedural shaders), and texture maps. As scene complexity increases, the amount of data required to render a scene on a local client increases as well. Assets for such complex virtual worlds are often included in a download or DVD that the user must first install before the 3D rendering engine on the local client can begin rendering a scene. Therefore, the more complex and dense the virtual world, the more data the user must either have installed or download on their system. Instantaneous virtual world rendering systems that depend on streaming render assets for client-side rendering may introduce objectionable visual artifacts as a viewport is procedurally refreshed with new textures and meshes that replace placeholder versions.\n4) Virtual world systems generally use a centralized server system to persistently maintain the state of the world. State data from these servers are continuously sent to each client when needed (including position, orientation, velocity and dynamics of objects within the world) in order for each viewport to remotely render the scene on the client's local machine. If all clients have the same assets, and the same hardware graphics capabilities, visual parity across each client viewport is possible.\nHowever, the overhead of synchronizing states across each user's viewport when the client systems are not homogenous can become problematic in many cases, as the visual representation of the same viewport may vary across each client, based on the client's local graphics capabilities. For example, if a faceted 3D object is tessellated and smoothed in hardware on a client system which supports curved patch tessellation, the silhouette of the object may appear to be significantly different than the silhouette of the same object on a client that does not support tessellation. In the latter case, the client would be forced to render the 3D object with a faceted outline. These differences can create objectionable discontinuities among remotely rendered viewports in cases where collision detection is performed by the server on the aforementioned 3D object. If the server performs collision detection using the smooth tessellated version of the object, and then transmits this singular result to each client, the visual representation of the point of impact and bounce vector from the collision may not appear to correspond to a valid point and slope on the surface of the object from the point of view of a user displaying the scene without tessellation. This can become even more objectionable, when the tessellation is performed on a macroscopic object such as terrain, where control of vehicles or avatars depends heavily on an accurate representation of the ground surface. These types of discontinuities between the state of world on the server, and the visual representation on an arbitrary client, increases significantly as a 3D virtual world becomes more complex, and the range of client devices that need to be supported expands to mobile phones, TV set top boxes and other hardware that may have limited or no 3D graphics capabilities on the local machine."}
{"text": "1. Field of the Invention\nThe invention relates to the field of information security and privacy and, more specifically, to finding matching entries in sets from two or more parties without revealing all of any party's entries.\n2. Discussion of the Prior Art\nModern research and technology have provided society with a wide variety of electronic devices. Some of these modern electronic devices are very powerful and useful to those who use them. For example, the global positioning system (GPS) was initially deployed for United States military purposes, but today can be a useful aid for almost anyone.\nFor back country travelers, boaters on open waters, or others who navigate off the beaten path, carrying a light weight, battery powered GPS receiver can provide them with their position within several hundred feet. For motorists, a dramatic improvement in convenience and safety can be provided by carrying a GPS receiver integrated with a trouble alert system (perhaps using mobile telephone technology) or a road map database.\nGPS receivers can be configured to record the time and position of the trajectories along which they travel. Such information can be maintained within such GPS receiver/recorder device. Periodically, it can be off-loaded onto a personal computer (PC), a personal digital assistant (PDA), or a server, for example.\nSuch trajectory information can be gathered and recorded with little or no intrusion on the person carrying the GPS receiver/recorder. Thus, a person could record their daily locations and travel over a period of days, years or even indefinitely.\nNevertheless, mining or making use of such information is difficult, both because of the volume of information potentially involved, and because of the privacy and security implications of sharing such information. Few individuals would be willing to make available to anyone else a complete log of where they have been over a period of time of significant duration. Even if certain “private” time periods could be easily deleted or not recorded, many people would consider a voluminous log of their whereabouts to be no business of anyone else.\nAdditionally, the security implications are daunting. Someone interested in burglarizing even an ordinary residence or place of business would be given a tremendous advantage by knowing the trends or patterns of who comes and goes from the place at what time—let alone if the person carrying the GPS receiver/recorder were a high profile person such as a leading politician or famous entertainer.\nOn the other hand, many people would be willing or even eager to reveal a limited subset of such trajectory information. For example, it often happens that two individuals who are meeting each other socially or for business have a feeling that they might have met before, but are unable to recall where and when that may have happened.\nMore generally, two or more individuals or organizations may possess information that can be represented as collections, data points or sets. They may desire to find what entities, elements or individual data points are present in all of the sets, but only if they can do so without revealing to the other participants all of the elements in their set. One example of such a problem includes several diplomats attempting to select where and when they will meet. Another example is two individuals determining if they both belong to any of the same golf courses, fraternal clubs or other organizations, which can be represented as the longitude and latitude of the clubhouse or headquarters.\nYet another example is two individuals selecting the time, date and place for initially meeting on a blind date. Another example is a company whose trucks are equipped with GPS receiver/recorder devices demonstrating that none of their trucks were near the coordinates of an accident. Yet another example is a person demonstrating that they have never gone near a location or region that is forbidden by the conditions of their parole or probation."}
{"text": "Non-volatile semiconductor memories (NVSMs) are widely used in many electronic devices such as personal digital assistants (PDAs), laptop computers, mobile phones and digital cameras. Some of these memories have arrays of charge storage transistors, such as floating gate transistors."}
{"text": "The present invention relates to a mechanism in a plastic injection machine and the mechanism ensures the screw rod to be precisely moved together with the open and close of the molds.\nA conventional plastic injection molding machine for manufacturing products having screw holes generally employs a gear system for controlling a screw rod moving into or removing from the cavity between two molds. The gear system is controlled by electric control system. The gear system including too many gears which is a load of maintenance. The gear system is difficult to precisely move the screw rod relative the molds so that the depth of the screw rod in the cavity cannot be obtained the same each time when the plastic material is injected into the cavity. The screw hole of the products could be damaged if the screw rod cannot be removed from the molds in stable. The gear system is activated after the movement of the molds is finished so that it prolongs the period of the complete process.\nIn accordance with one aspect of the present invention, there is provided a plastic injection molding machine which comprises a position block on the upper mold and a slide member on the lower mold. A pin is inserted in the position block and the movable member. A rod extends through the slide member so that the top end of the rod is inserted in the cavity when the pin is inserted in the slide member. Two side members each have threaded semi-circular recesses which are engaged with a threaded section on the rod so that when the rod is rotated by a motor, the rod is smoothly removed from the cavity.\nThe present invention will become more obvious from the following description when taken in connection with the accompanying drawings which show, for purposes of illustration only, a preferred embodiment in accordance with the present invention."}
{"text": "This application is based upon and claims the benefit of priority from the prior Japanese Patent Applications No. 2000-117739, Apr. 19, 2000; No. 2000-117740, Apr. 19, 2000; No. 2001-055815, Feb. 28, 2001; and No. 2001-110469 Apr. 9, 2001, the entire contents of which are incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to a focus stabilizing apparatus for stabilizing a focused state of an optical apparatus, such as a microscope.\n2. Description of the Background Art\nGenerally, a sample observation is carried out with the use of a microscope. This sample observation is made as follows. An observation sample is placed on a microscope stage and an objective lens is moved closer to the observation sample. By doing so, an observation spot of the observation sample is observed under a magnified state.\nIn this case, an objective lens has its focal depth decreased as the magnifying power becomes higher. It is, therefore, difficult to achieve a focus setting between the objective lens and the observation sample. Further, if the distance between the objective lens and the observation sample minutely varies, defocusing occurs between the focal point position of the objective lens and the observation sample due to a variation of the minute distance, so that the quality of the observation image is greatly degraded.\nOn the other hand, an apparent position between the objective lens and the observation sample is very closer to each other. By the way, there exists a mechanical coupling length between the objective lens and the observation sample. This mechanical coupling length is constituted by many mechanical component parts present between the objective lens and the observation sample. The mechanical coupling length is provided by the length over which, for example, a microscope frame, objective lens moving mechanism, objective lens mounting revolver are passed. Therefore, the length is very long because many mechanical components are interposed.\nThese mechanical component parts are liable to be varied in their dimensions due to a temperature variation involved.\nFurther, the greater the number of the mechanical component parts the greater the mechanical coupling length involved. The microscope is easily affected by a vibration and the vibration amplitude becomes greater.\nIf, therefore, the ambient temperature varies due to, for example, the turning ON/OFF of an illumination, internal power supply, etc., as well as the operation of an air conditioning equipment, then there arises a variation in the dimensions of mechanical component parts in the microscope. Even if, therefore, the focal setting of the objective is made relative to the observation sample, the distance between the objective lens and the observation sample greatly varies due to the above-mentioned dimensional variation and there occurs a defocusing.\nFurther, under a somewhat smaller external vibration, a greater vibration amplitude is involved and a distance between the objective lens and the sample varies, thus resulting in an out-of-focus state.\nHeretofore, various kinds of autofocusing mechanisms have been considered to compensate such defocusing. These autofocusing mechanisms require a complex mechanical/electrical mechanism and control system. For this reason, the resultant apparatus becomes bulkier and expensive. Further, another microscope is known in which, like a fluorescent observation, the brightness is extremely darker at a time of observation. To such a microscope it is originally difficult to apply the above-mentioned autofocusing mechanism.\nThe following technique is disclosed in JPN PAT APPLN KOKAI Publication No. 9-120030. In this technique, a stage is provided through a rack and drive gears and it is driven in the optical axis direction of the objective lens. Between the gears and the stage, at least two rods are inserted. These rods are different in their thermal expansion coefficients and configured such that the direction of the thermal expansion coefficient of one rod acts in a direction opposite to that of the other rod. By doing so, defocusing is compensated.\nIn such a technique, the rods are arranged within the microscope body and a time is taken until the temperature is compensated under a variation of the ambient temperature. For this reason, there is a risk that, when the sample is observed, its operation, efficiency will be lowered. Further, since the rods are interposed, the mechanical coupling length becomes greater and its structure is liable to be affected from the external vibration. It is, therefore, necessary to remodel the microscope itself.\nIn the case where a living cell, etc., is observed as a target object, only a periphery side of the observation sample is sometimes warmed by a warmer. However, it is not possible to compensate a temperature drift involved.\nIt is accordingly the object of the present invention to provide a focus stabilizing apparatus which can make a focus setting of an objective lens relative to an observation sample under a stable way without being adversely influenced from a variation in the ambient temperature.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to a observation sample, and position adjusting means provided at an outer periphery of the objective lens and adapted to move the observation sample in an optical axis direction to set a focal point of the objective lens to the observation sample.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to the observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and adapted to move the observation sample in an optical axis direction through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set through the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an objection lens, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and adapted to move the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base for placing the observation sample thereon, and an operation ring threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens to move the observation sample in the optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base placing the observation sample thereon, and an operation ring threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the moving axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on a stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, and a sample retaining base placing the observation sample thereon and threadably inserted over the fixing base and, being rotated relative to the fixing base, moving in the optical axis direction of the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens and a sample retaining base placing the observation sample thereon and threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens.\nIn such focus stabilizing apparatus, the objective lens is fixed to the fixing base and threadably inserted directly into a revolver provided at a microscope body.\nThe position adjusting means has a rotation stop section arranged at the upper side of the sample retaining base, an intermediate seat arranged on the upper side of the rotation stop section and having the observation sample placed thereon, and a pin having one end mounted in the fixing base and the other end inserted through a hole in the rotation stop section and provided parallel to the optical axis direction of the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base placing the observation sample thereon, fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, and a friction member situated between the fixing base and the sample retaining base and retaining the sample retaining base by a friction force relative to the fixing base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain an observation sample, position adjusting means provided on an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base placing the observation sample thereon and fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, and a friction member situated between the fixing base and the sample retaining base and retaining the sample retaining base by a friction force relative to the fixing base.\nIn the focus retaining apparatus, a sample retaining base grasping mechanism is provided for grasping and fixing the sample retaining base.\nFurther, a temperature sensor is provided on the fixing base to detect the ambient temperature and a temperature adjusting means is provided for adjusting the temperature of the fixing base and sample retaining base on the basis of a detection output of the temperature sensor to be made constant at all times.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, a piezoelectric actuator provided on the upper side of the sample retaining base and performing its extending/contracting operation in the optical axis direction of the objective lens, a sample base provided on the upper side of the piezoelectric actuator and placing the observation sample thereon and an electrostatic sensor provided on the upper side of the sample retaining base to detect a moving amount of the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens, a sample retaining base fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens, a piezoelectric actuator provided on the upper side of the sample retaining base and performing its extending/contracting operation in the optical axis direction of the objective lens, a sample base provided on the upper side of the piezoelectric actuator and placing the observation sample thereon, and an electrostatic sensor provided on the upper side of the sample retaining base to detect a moving amount of the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, and position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, wherein the position adjusting means has a fixing base for fixing the objective lens and having a flange at its outer periphery and a sample retaining base fitted over an outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens and having a flange at its outer periphery, and a feed screw section provided between the flange of the fixing base and the flange of the sample retaining base to feed the sample retaining base in the optical axis direction of the objective lens relative to the fixing base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation stage, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage, wherein the position adjusting means has a fixing base for fixing the objective lens and having a flange at its outer periphery, a sample retaining base fitted over the outer periphery of the fixing base and so provided as to be movable in the optical axis direction of the objective lens and having a flange at its outer periphery, and a feed screw section provided between the flange of the fixing base and the flange of the sample retaining base to feed the sample retaining base in the optical axis direction of the objective lens relative to the fixing base.\nIn the focus stabilizing apparatus, leaf springs are provided on the sample retaining base to fix the observation sample thereon.\nFurther, a mechanical coupling length between the objective lens and observation sample is set by the fixing base, sample retaining base and operation ring.\nA mechanical coupling length between the objective lens and the observation sample is set by the fixing base and the sample retaining base.\nIn the focus stabilizing apparatus, the fixing base, operation ring and sample retaining base are made of at least two different kinds of materials and selectable in their dimensions.\nThe operation ring and sample retaining base are made of at least two different kinds of materials and selectable in their dimensions.\nThe operation ring and sample retaining base are formed of materials of different linear expansion coefficients.\nAn auxiliary member is located between the operation ring and the sample retaining base and made of a material different in liner expansion coefficient from the materials of the operation ring and sample retaining base.\nThe elastic means is comprised of leaf springs provided on the stage and arranged along the direction orthogonal to the optical axis of the objective lens and a magnet provided on the leaf spring to attract the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens to set a focal point of the objective lens to the observation sample, a minute movement mechanism for minutely displacing the objective lens in the optical axis direction of the objective lens, displacement amount detecting means for detecting a displacement amount of the objective lens, and control means for operating the minute movement mechanism on the basis of a detection output of the displacement amount detecting means to control a relative distance between the objective lens and the observation sample to a predetermined distance.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a stage so provided as to be movable in a direction orthogonal to an optical axis of the objective lens, a sample base placed on the stage to retain the observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens through the sample base to set a focal point of the objective lens to the observation sample, a minute movement mechanism for minutely displacing the objective lens in the optical axis direction of the objective lens, displacement amount detecting means for detecting a displacement amount of the objective lens, control means for operating the minute movement mechanism on the basis of a detection output of the displacement amount detecting means to control a relative distance between the objective lens and the observation sample to a predetermined distance, and elastic means provided on the stage and so set relative to the sample base as to have a low rigidity in the optical axis direction of the objective lens and a high rigidity in a moving direction of the stage.\nIn the focus stabilizing apparatus, the position adjusting means has a fixing base for fixing the objective lens through the minute movement mechanism, a sample retaining base for placing the observation sample thereon, and an operation ring threadably inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nA mechanical coupling length between the objective lens and the observation sample is set by the fixing base for fixing the objective lens through the minute movement mechanism, the sample retaining base for placing the observation sample thereon and the operation ring inserted over the fixing base and, by being rotated relative to the fixing base, moving in the optical axis direction of the objective lens and moving the sample retaining base in the optical axis direction of the objective lens to adjust a relative distance between the observation sample and the objective lens.\nThe minute movement mechanism has a moving stage relative to which the objective lens is provided, and an actuator for minutely moving the moving stage in the optical axis direction of the objective lens.\nThe minute movement mechanism has a moving stage relative to which the objective lens is provided, and piezoelectric actuators arranged in those positions symmetric relative to the optical axis of the objective lens to minutely move the moving stage in the optical axis direction of the objective lens.\nThe displacement amount detecting means is comprised of an electrostatic capacity sensor.\nThe control means has means for receiving, as inputs, an objective lens displacement amount detected by the displacement amount detecting means and an instruction value representing the position of the objective lens, means for finding a deviation between the displacement amount of the objective lens and the instruction value, and means for performing operation control of the minute movement mechanism in accordance with the deviation to move the objective lens to a position designated by the instruction value.\nThe elastic means is comprised of leaf springs mounted on the stage and arranged along the direction orthogonal to the optical axis of the objective lens and a magnet mounted on the leaf spring to attract the sample base.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a minute movement mechanism for minutely displacing the objective lens in the optical axis direction of the objective lens, displacement amount detecting means for detecting a displacement amount of the objective lens, and control means for operating the minute movement mechanism on the basis of a detection output of the displacement amount detecting means to control a relative distance between the objective lens and the observation sample to a predetermined distance.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, a fixing base for fixing the objective lens, a minute movement table on which the objective lens is provided, parallel springs situated between the fixing base and the minute movement table to allow the minute movement table to be moved in an optical axis direction of the objective lens, an actuator provided between the fixing base and the minute movement table to minutely displace the minute movement table in the optical axis direction of the objective lens, a displacement sensor for detecting a displacement amount of the objective lens, and control means for allowing the actuator to perform its extending/contracting operation on the basis of a detection output of the displacement sensor to control the objective lens and bring it to a just-in-focus position relative to the observation sample.\nIn the focus stabilizing apparatus, the control means has a memory section for storing an output of the displacement sensor corresponding to a just-in-focus state between the observation sample and the objective lens, a comparing section for comparing an output of the displacement sensor and an output of the displacement sensor stored in the memory section, and a control section for outputting an electric signal for canceling a distance variation between the observation sample and the objective lens on the basis of a result of comparison by the comparing section to the actuator.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises an objective lens arranged opposite to an observation sample, position adjusting means provided at an outer periphery of the objective lens and moving the observation sample in an optical axis direction of the objective lens to set a focal point of the objective lens to the observation sample, a minute movement mechanism for minutely varying a distance between the observation sample and the objective lens in the optical axis direction, and control means for operating the minute movement mechanism to make a fine adjustment of the distance between the observation sample and the objective lens.\nAccording to a main aspect of the present invention, a focus stabilizing apparatus is provided which comprises a fixing base having an objective lens mounting section, a sample base for supporting an observation sample, and position adjusting means provided on the fixing base to allow the sample base to be moved in an optical axis direction of the objective lens.\nIn the focus stabilizing apparatus, the position adjusting means has a fixing base for fixing the objective lens, an operation ring inserted over the fixing base and, by being rotated, moving in the optical axis direction of the objective lens, and a sample retaining base so provided as to be movable in the optical axis direction and, upon a rotation operation of the operating ring, adjusting a distance between the observation sample and the objective lens.\nAccording to the present invention, the mechanical coupling length between the objective lens and the observation sample is determined by only the position adjusting means provided relative to the objective lens and can be set to be very short. By doing so, there is almost no variation in a positional relation between the objective lens and the observation sample even if the ambient temperature varies. The structure suffers no adverse effect from an external vibration and the sample can be observed at all times under a stable condition.\nAfter the focal point of the objective lens has been set to the observation sample, the observation sample can be moved in a direction orthogonal to the optical axis of the objective lens simply by moving the stage and it is also possible to readily position the observation sample in a moving direction of the stage.\nFurther, as the position adjusting means use is made of at least two different kinds of materials and its component elements are selectable in their materials and dimensions. It is, therefore, possible to prevent an observation image from being degraded even if the ambient temperature varies.\nFurther, the mechanical coupling length between the objective lens and the observation sample is set to be short and the operation of the minute movement mechanism is controlled in accordance with a deviation between the displacement amount of the objective lens and the instruction value to allow the objective lens to be set to a position designated by the instruction value. Therefore, even if the ambient temperature varies, there is almost no variation in positional relation between the objective lens and the observation sample. The structure also suffers no adverse effect from external vibrations and the objective lens can be set to a desired position given by the instruction value and the observation sample can be observed, under a stable way, over an extended period of time.\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out hereinafter."}
{"text": "1. Field of the Invention\nAn aspect of the present invention relates to a thin film deposition apparatus and a method of manufacturing an organic light emitting device (OLED) by using the same, and more particularly, to a thin film deposition apparatus to form a fine pattern on a large-sized substrate and a method of manufacturing an OLED by using the same.\n2. Description of the Related Art\nA display device using an organic light emitting device (OLED) has a wide viewing angle, excellent contrast, and rapid response time and thus has recently been highlighted as a next generation display device.\nIn general, an OLED emits light when holes and electrons respectively injected from an anode and a cathode are recombined in an organic light emitting layer, thereby generating visible light. However, in this structure, it is difficult to obtain high light emission efficiency. Thus, in an OLED, other organic layers such as an electron injection layer, an electron transport layer, a hole transport layer, and a hole injection layer may be selectively disposed between the anode and the cathode along with the organic light emitting layer.\nMoreover, electrodes, the organic light emitting layer, and the organic layer of the OLED may be formed using various methods including deposition. In order to manufacture an organic light emitting display device using deposition, a fine metal mask (FMM) having a pattern of a thin film to be formed is closely adhered to the surface of a substrate on which the thin film is to be formed, and a material for forming the thin film is deposited on the FMM, thereby forming a thin film having a predetermined pattern.\nHowever, it is difficult to form a fine pattern of organic thin films and red, green, and blue luminous efficiency may vary according to the pattern and thickness of the organic thin films. Thus, there are limitations in improving luminous efficiency of an OLED.\nIn addition, while the size of current display devices is gradually increasing, it is difficult to pattern an organic thin film over a large area by using a conventional thin film deposition apparatus. Thus, it is difficult to manufacture a large-sized OLED having satisfactory driving voltage, current density, brightness, color purity, luminous efficiency, and life span."}
{"text": "This invention relates to high temperature electrochemical converters and specifically to high performance systems employing such devices and methods.\nElectrochemical converters perform fuel-to-electricity conversions in a fuel cell (electric generator) mode or electricity-to-fuel conversions in an electrolyzer (fuel synthesizer) mode. The converters are capable of high efficiencies, depending only on the relation between the free energy and enthalpy of the electrochemical reaction, and are not limited by Carnot-cycle considerations.\nThe key components in an electrochemical energy converter are a series of electrolyte units onto which electrodes are applied and a similar series of interconnectors disposed between the electrolyte units to provide serial electrical connections. Each electrolyte unit is an ionic conductor having low ionic resistance thereby allowing the transport of an ionic species from one electrode-electrolyte interface to the opposite electrode-electrolyte interface under the operating conditions of the converter. Various electrolytes can be used in such converters. For example, zirconia stabilized with such compounds as magnesia, calcia or yttria can satisfy these requirements when operating at an elevated temperature (typically around 1000.degree. C.). The electrolyte material utilizes oxygen ions to carry electrical current. The electrolyte should not be conductive to electrons which can cause a short-circuit of the converter. On the other hand, the interconnector must be a good electronic conductor. The interaction of the reacting gas, electrode and electrolyte occurs at the electrode-electrolyte interface which requires the electrodes be sufficiently porous to admit the reacting gas species and to permit exit of product species.\nThe approach of forming electrolyte and interconnector components as free-standing plates was disclosed by the present inventor in U.S. Pat. No. 4,490,445, issued Dec. 25, 1984, which is herein incorporated by reference. However, during operation, the stacks of electrolyte and interconnector plates can experience thermal non-equilibrium. Thus, it is important to reduce thermal gradients across the entire converter assembly by facilitating the transfer of heat to and from the electrolyte elements.\nWhen an electrochemical converter performs fuel-to-electricity conversion in a fuel cell mode, waste energy in the form of heat must be removed from the electrolyte surfaces. Conversely, when the converter performs electricity-to fuel conversion in the electrolyzer mode, the electrolyte must be provided with heat to maintain its reaction. In prior systems, heat exchanging has been achieved primarily by the convective heat transfer capabilities of the gaseous reactants as they travel through the assembly. Such reliance on the heat capacity of the reactants creates an inherent thermal gradient in the system, resulting in non-optimum electrochemical processes.\nTo rectify this problem, the approach of integrating with the electrochemical converter a series of heat transfer elements was disclosed by the present inventor in U.S. Pat. No. 4,853,100, issued Aug. 1, 1989, which is herein incorporated by reference. The above-mentioned integration system facilitates the heat transfer from the fuel cell stacks by reducing thermal gradients across the converter assembly. However, there still exists a need for further improvements in the thermal control mechanisms within electrochemical energy systems. In particular, an improved electrochemical energy conversion system having the ability to more efficiently regulate the operating temperature within the electrochemical assembly would represent a major improvement in the industry."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a server computer system that provides centralized computing resources to a plurality of users. More particularly, the invention relates to a system and method for maintaining a pool of free virtual machines available to be assigned to users on a server computer.\n2. Description of the Related Art\nMany commercial businesses and enterprises make extensive use of personal computers (PCs) in their daily operations. Typically, each user of a PC in the enterprise has a networked PC at his/her desk or work area. As the number of networked computer systems utilized in an enterprise increases, the management of resources in the network may become increasingly complex and expensive. Some of the manageability issues involved in maintaining a large number of networked computer systems may include ease of installation and deployment, the topology and physical logistics of the network, asset management, scalability (the cost and effort involved in increasing the number of units), troubleshooting network or unit problems, support costs, software tracking and management, as well as the simple issue of physical space, be it floor space or room on the desktop, as well as security issues regarding physical assets, information protection, software control, and computer virus issues.\nMany of these issues may be addressed by centralizing the locations of computing resources. For example, each individual user may connect through a network to a server computer system and use computing resources provided by the server computer system.\nSome server computer systems are capable of leveraging the physical hardware resources available by providing virtualization capabilities. Virtualization provides the ability for multiple virtual machines to run together on the same physical server computer. For example, each virtual machine may execute its own operating system and may appear to a user of the virtual machine to be the same as an independent physical computer. The software layer that executes on the physical server computer and manages the various virtual machines is called a hypervisor or virtual machine host software. The hypervisor can run on bare hardware (called a Type 1 or native VM) or under control of an operating system (called a Type 2 or hosted VM)."}
{"text": "1. Field of the Invention\nThe present invention relates to an inkjet recording apparatus which ejects ink droplets to record an image on a recording medium.\n2. Description of the Related Art\nJapanese Unexamined Patent Publication No. 59597/2005 (Tokukai 2005-59597) discloses an inkjet printer including a conveyance mechanism having a drum which rotates to convey a sheet carried on the outer circumferential surface thereof; a plurality of inkjet heads each having an ejection surface, which are aligned in a conveyance direction of the sheet so that the ejection surface of the each inkjet heads faces the outer circumferential surface of the drum; and a wiper for wiping the ejection surface. In this inkjet printer, all the inkjet heads are fixed on a frame structure. The frame structure is moveable between a printing position and a wiping position. The printing position is a position where the frame structure is disposed when ink droplets are ejected from an ejection surface to a sheet conveyed by the conveyance mechanism. The wiping position is such a position that the ejection surface is disposed farther apart from the outer circumferential surface of the drum, compared to the printing position. At a time of printing, the wiper is in a standby position and faces no ejection surfaces, and the frame structure is positioned in the printing position. At a time of a wiping operation, the frame structure moves to the wiping position. Then, the wiper moves in the circumferential direction of the drum, from the standby position to an opposing position so as to face the ejection surface. Then, the wiper reciprocates in the axial direction of the drum, thereby wiping the ejection surface."}
{"text": "Wireless communication systems are widely deployed to provide various types of communication content such as voice, data, and so on. These systems may be multiple-access systems capable of supporting communication with multiple users by sharing the available system resources (e.g., bandwidth and transmit power). Examples of such multiple access systems include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, 3GPP Long Term Evolution (LTE) systems, and orthogonal frequency division multiple access (OFDMA) systems.\nGenerally, a wireless multiple-access communication system can simultaneously support communication for multiple wireless terminals. Each terminal communicates with one or more base stations via transmissions on the forward and reverse links. The forward link (or downlink) refers to the communication link from the base stations to the terminals, and the reverse link (or uplink) refers to the communication link from the terminals to the base stations. This communication link may be established via a single-in-single-out, multiple-in-single-out or a multiple-in-multiple out (MIMO) system.\nSome conventional advanced devices include multiple radios for transmitting/receiving using different Radio Access Technologies (RATs). Examples of RATs include, e.g., Universal Mobile Telecommunications System (UMTS), Global System for Mobile Communications (GSM), CDMA2000, WiMAX, WLAN (e.g., WiFi), Bluetooth, LTE, and the like.\nAn example mobile device is an LTE User Equipment (UE), such as a fourth generation (4G) mobile phone. Such 4G phone may include various radios to provide a variety of functions for the user. For purposes of this example, the 4G phone includes an LTE radio for voice and data, an IEEE 802.11 (WiFi) radio, a Global Positioning System (GPS) radio, and a Bluetooth radio, where two of the above or all four may operate simultaneously. While the different radios provide useful functionalities for the phone, their inclusion in a single device gives rise to coexistence issues. Specifically, operation of one radio may in some cases interfere with operation of another radio through radiative, conductive, resource collision, and/or other interference mechanisms. Coexistence issues include such interference.\nThis is especially true for the LTE uplink channel, which is adjacent to the Industrial Scientific and Medical (ISM) band and may cause interference therewith It is noted that Bluetooth and some Wireless LAN (WLAN) channels fall within the ISM band. In some instances, a Bluetooth error rate can become unacceptable when LTE is active in some channels of Band 7 or even Band 40 for some Bluetooth channel conditions. Even though there is no significant degradation to LTE, simultaneous operation with Bluetooth can result in disruption in voice services terminating in a Bluetooth headset. Such disruption may be unacceptable to the consumer. A similar issue exists when LTE transmissions interfere with GPS. Currently, there is no mechanism that can solve this issue because LTE by itself does not experience any degradation\nWith reference specifically to LTE, it is noted that a UE communicates with an evolved NodeB (eNB; e.g., a base station for a wireless communications network) to inform the eNB of interference seen by the UE on the downlink. Furthermore, the eNB may be able to estimate interference at the UE using a downlink error rate. In some instances, the eNB and the UE can cooperate to find a solution that reduces interference at the UE, even interference due to radios within the UE itself. However, in conventional LTE, the interference estimates regarding the downlink may not be adequate to comprehensively address interference.\nIn one instance, an LTE uplink signal interferes with a Bluetooth signal or wireless local area network (WLAN) signal. However, such interference is not reflected in the downlink measurement reports at the eNB. As a result, unilateral action on the part of the UE (e.g., moving the uplink signal to a different channel) may be thwarted by the eNB, which is not aware of the uplink coexistence issue and seeks to undo the unilateral action. For instance, even if the UE re-establishes the connection on a different frequency channel, the network can still handover the UE back to the original frequency channel that was corrupted by the in-device interference. This is a likely scenario because the desired signal strength on the corrupted channel may sometimes be higher as reflected in the measurement reports of the new channel based on Reference Signal Received Power (RSRP) to the eNB. Hence, a ping-pong effect of being transferred back and forth between the corrupted channel and the desired channel can happen if the eNB uses RSRP reports for handover decisions.\nOther unilateral action on the part of the UE, such as simply stopping uplink communications without coordination of the eNB may cause power loop malfunctions at the eNB. Additional issues that exist in conventional LTE include a general lack of ability on the part of the UE to suggest desired configurations as an alternative to configurations that have coexistence issues. For at least these reasons, uplink coexistence issues at the UE may remain unresolved for a long time period, degrading performance and efficiency for other radios of the UE."}
{"text": "1. Field of the Invention\nThe present invention is directed to a ground contact height sensor for a harvesting header having two pivotally coupled rigid arms.\n2. Description of the Prior Art\nAgricultural combines may be provided with header height control systems to control the position of the harvesting header relative to the ground. Ground contact sensors may be used to determine the distance of the bottom of the header to the ground. These ground contact sensors generally comprise a downwardly extending arm that contacts the ground below the header. As the arm is deflected upwardly by an obstruction the arm moves a mechanical linkage or a potentiometer which signals the automatic header height control system raising the harvesting header.\nIn one sensor configuration a J-shaped rigid arm is coupled to a transverse rotatable shaft. As the arm is deflected the shaft is rotated shifting a potentiometer or a mechanical linkage. Because of the J-shape and the rigid nature of the arm it can be difficult to store these arms inside the floor of the header. In addition, these J-shaped arms have a relatively narrow radius, so it can be difficult to put the combine in reverse without damaging the sensor.\nIn another sensor configuration bowed spring steel hoops form the arms. The hoops can be flattened out to fit within a recess under the floor of the platform. In addition, the hoops form a large radius facilitating reverse movement of the combine. However hoop sensors can be easily damaged during sharp turns if they are following a furrow.\nIt is an object of the present invention to provide a ground contact height sensor that can be easily stored within a recess under the floor of the header, that will not be easily damaged during sharp turns, and that will not be easily damaged during reverse movements.\nThe sensor comprises a transverse rotatable shaft, a first rigid arm, a second rigid arm and a connecting member. The rotatable shaft is rotatively mounted to the frame of the header. The first rigid arm has a first end that is rigidly mounted to the rotatable shaft. The first arm extends downwardly and rearwardly from the rotatable shaft in a curved manner. The second rigid arm is pivotally coupled to the second end of the first rigid arm. The second rigid arm extends upwardly and rearwardly from the second end of the first rigid arm. The first and second rigid arms defining a relatively large radius. A connecting member secures the rearmost end of the second rigid arm to the frame of the header. In the illustrated embodiment the connecting member comprises a chain. To maintain the radius of the two pivotally connected rigid arms the hinge joint joining the two rigid arms is provided with a torsion spring. The limits of the pivoting movement is controlled by stops."}
{"text": "This invention relates to a black pigment, a process for making such a pigment and paints and/or plastics containing this pigment.\nHeat stable black pigments which contain no carbon have generally been made from materials containing heavy metals and/or other highly toxic materials. Such materials are described, for example, in the Colour Index published in Great Britain and the brochure entitled \"Classification and Chemical Descriptions of Mixed Metal Oxide Inorganic Colored Pigments\" published by the Dry Color Manufacturers' Association, 1117 North 19th Street, Suite 100, Arlington, Va., 22209. These heavy metal-containing and/or highly toxic materials are, however, undesirable in many applications."}
{"text": "1. Technical Field.\nThis invention pertains to fluid control systems and particularly to valves. The apparatus of this invention is useful for preventing and correcting the deleterious affects of marine mussels colonizing in marine raw water inlet systems of boats and equipment. The valve apparatus is useful for winterizing these systems where introduction of antifreeze compounds may be desired.\n2. Background Information.\nRaw water inlet systems, such as those commonly employed in boat engines, marine heads, and marine generators, operate most efficiently when a free and unrestricted flow of water is able to pass through them. The water is used for cooling of engines or machinery, and for flushing of sanitation devices. If the flow of water should become restricted or stopped, the effected equipment may be disabled or damaged.\nIt has been determined that some marine animals, particularly zebra mussels, dreissena polymorpha, will colonize in the raw water inlets of marine equipment thereby disabling or damaging the equipment. It is therefore economically advantageous to control or eliminate these colonies in boats and marine equipment. The zebra mussel is a small bivalve mollusk which has an elongated thick shell showing alternating light and dark bands. The zebra mussel is found in all of the Great Lakes and is expected to proliferate into other freshwater systems in the Midwestern and Eastern portions of the United States. The zebra mussel is believed to have a substantial negative impact on commercial and recreational boating, including damage to boat cooling and sanitation systems due to constriction of boat plumbing by colonies of the animal. Current zebra mussel control techniques involve scraping or application or chemicals such as hydrogen sulfide, chlorine, ozone, and copper sulfate. It is known that zebra mussels are sensitive to temperatures above 37 degrees C. (98.6.degree. F.), are quickly killed at temperatures between 45.degree.-55.degree. C. (113.degree.-131.degree. F.) and are immediately killed at temperatures above 60.degree. C. (140.degree. F.). This property, sensitivity to heat and chemicals, is utilized by the apparatus and method of the present invention to control zebra mussel infestation.\nIn addition to the problems associated with zebra mussel infestation in marine plumbing systems, boats and equipment that are laid up in winter due to freezing conditions are subject to damage of their machinery and plumbing systems from internal build up of ice. It is therefore necessary to winterize this equipment by draining and/or the introduction of antifreeze compounds into the raw water systems of the equipment.\nIn the past, various devices and/or methods have been used or proposed, to either control zebra mussel colonization or permit marine system winterization. However, these devices and methods have significant limitations and shortcomings. And no single apparatus or method is known to have been used for both purposes.\nDespite the need in the art for an apparatus and method which overcome the shortcomings and limitations of the prior art, none insofar as is known has been developed or proposed. Accordingly, it is an object of the present invention to provide an apparatus and method of use therefor, which prevents and corrects zebra mussel colonization in marine raw water inlet systems of boats and other equipment. It is also an object of this invention to provide an apparatus and method which is usable to introduce antifreeze and other compounds into the raw water inlet systems of boats and other equipment for winterization purposes. It is an additional object of the invention to provide a combined mussel control and antifreeze introduction system using a single, unitary fittng or set of fittings disposed in the marine equipment. It is a further object of this invention to provide a Zebra mussel control and marine winterization apparatus which is reliable, efficient, economical, and usable by operators and service personnel at marinas and boat yards, as well as owner operators of marine equipment, and which overcomes the limitations and shortcomings of the prior art."}
{"text": "This invention relates to a method and control for defrosting the outdoor coil of a heat pump in a manner which optimizes efficiency and conserves energy.\nWhen a heat pump operates in its heating mode, frost builds up on the pump's outdoor coil. As the frost thickness increases, heat transfer from the outdoor air decreases and the efficiency of the heat pump drops significantly, a substantial amount of energy therefore being wasted. Hence, it is necessary to periodically defrost the outdoor coil. This is usually accomplished by reversing the refrigerant flow in the heat pump which will heat the outdoor coil and melt the frost.\nIt is recognized that there is an optimum point of frost accumulation at which the heat pump should be switched to its defrost mode of operation. If defrost is commanded too soon or too late, energy will be wasted and efficiency will suffer. Unfortunately, it has been very difficult to achieve such optimum operation in the past. Moreover, these previous defrost systems are unreliable in operation and/or are not adaptable to all types of outdoor coils.\nSubstantially less expensive defrost control systems have also been developed, but these systems are not capable of adjusting to the prevailing weather conditions. In one such system, the differential between the outdoor ambient (dry bulb) temperature and the refrigerant temperature in the outdoor coil is measured. The outdoor coil temperature decreases as frost builds up, and this increases the temperature split or difference between the outdoor ambient temperature and the coil temperature. When the temperature split increases to a predetermined value, the outdoor coil is defrosted. These prior temperature differential type defrost controls, however, fail to take the weather conditions into account. The temperature split between the outdoor ambient air (dry bulb) temperature and the refrigerant temperature in the outdoor coil for clean coil operation is a function of the outdoor wet bulb temperature and not the dry bulb temperature. For example, when the outdoor ambient air has a 35.degree. F. dry bulb temperature, a 34.degree. F. wet bulb temperature, and a relative humidity of about 90%, the refrigerant temperature in the outdoor coil of a typical three ton heat pump may be about 23.degree. F. when the outdoor coil is frost-free, the clean coil temperature split (namely, the outdoor ambient temperature minus the outdoor coil temperature) thereby being 35.degree.-23.degree. or 12.degree.. (All temperatures mentioned herein will be F or Fahrenheit.) For the same outdoor dry bulb temperature, an outdoor wet bulb temperature of 28.degree. and an outdoor relative humidity of about 40% may then provide an outdoor coil temperature of about 17.degree., resulting in a clean coil temperature split of 35.degree.-17.degree. or 18.degree.. Neither humidity condition is uncommon in most areas. Thus, if the defrost control were set, when the ambient air has a 34.degree. wet bulb temperature, to initiate defrost at a temperature differential of, for example, 5.degree. above its expected clean coil condition, defrost would occur when the temperature differential became 12.degree.+5.degree. or 17.degree. and dry weather conditions would result in the system continually defrosting itself without time for frost buildup on the outdoor coil.\nEven if the temperature split, at which defrost should occur, is properly determined when the outdoor coil is frost-free, long before frost builds up and that temperature split is reached the weather conditions (namely, the outdoor temperature and/or relative humidity) may change significantly, and that previously determined temperature split may no longer be appropriate or valid. If there is a decrease in outdoor temperature between defrost modes, excessive frost would build up on the outdoor coil and defrost should now be initiated at a smaller temperature split, not the one previously determined. On the other hand, as the outdoor temperature rises the same system may go into needless defrost because the control would assume that frost is building up on the coil, when it may not.\nThis phenomenon may be appreciated and more fully understood by observing FIG. 1 which provides a graph of the performance of the typical three ton heat pump mentioned previously. The graph plots the wet bulb temperature of the outdoor air versus the outdoor ambient or dry bulb temperature at different outdoor relative humidities. The graph shows the liquid line temperature, which is essentially the same as the outdoor coil temperature or the coil surface temperature, under clean coil conditions at various wet bulb temperatures. The clean coil temperature splits (the outdoor dry bulb temperature minus the liquid line temperature) for different weather conditions, namely at different points on the graph, may easily be determined by subtraction of one temperature from the other at the point that represents the weather conditions. The graph clearly illustrates that the liquid line temperature is strictly a function of the wet bulb temperature, and thus the moisture in the outdoor air.\nIt will be assumed that on a given day at about 7 a.m. the weather conditions in a particular area are as depicted by point 11 in FIG. 1, namely about 12.degree. outdoor ambient temperature, 10.5.degree. wet bulb temperature and about 77% relative humidity, the liquid line temperature for clean coil conditions thus being about 4.5.degree. to provide a clean coil temperature split of 12.degree.-4.5.degree. or 7.5.degree.. Point 12 indicates the assumed weather conditions on the same day at 10 a.m.--29.degree. outdoor dry bulb temperature, 23.degree. wet bulb temperature, about 40% relative humidity and a liquid line temperature of about 13.5.degree., the clean coil temperature split thereby being 29.degree.-13.5.degree. or 15.5.degree.. This corresponds to an 8.degree. increase (15.5-7.5) in the temperature split for a clean outdoor coil. If the control system were programmed, in accordance with the data at 7 a.m., to initiate defrost after there is a 4.degree. temperature increase in the clean coil temperature split, a needless defrost cycle would occur with no frost build up on the outdoor coil. Points 13 and 14 in FIG. 1 depict the assumed weather conditions at 4 p.m. and 11 p.m., respectively, on the same given day. The graph indicates that the clean coil temperature split would change downward from about 18.degree. to 11.5.degree., or about 6.5.degree., between 4 p.m. and 11 p.m. Thus, a 4.degree. programmed differential would require that the initial 18.degree. clean coil split at 4 p.m. would have to increase to 22.degree. before defrost would occur, whereas the optimum defrost split (the difference between the outdoor temperature and the coil temperature when the defrost mode should be initiated) for the weather conditions at 11 p.m. would be 11.5.degree. plus 4.degree., or 15.5.degree.. Hence, the split would increase 6.5.degree. (from 15.5.degree. to 22.degree.) above the optimum defrost condition before defrost would be initiated and excessive frost would accumulate. The conditions assumed in explaining the FIG. 1 graph are not uncommon, since the outdoor temperature and relative humidity may experience wide variations over a 24-hour period.\nThe defrost control system of the present invention is a substantial improvement over those previously developed. The system is not only relatively inexpensive but the initiation of outdoor coil defrost is timed to occur at the optimum point regardless of changing weather conditions so that defrost only and always occurs when it is necessary, thereby increasing the efficiency of the heat pump, conserving energy and improving system reliability. Any time there is a significant change in the weather conditions, the control system of the present invention will effectively recalculate when a defrost cycle should be initiated."}
{"text": "Compounds of Formula I: ##STR1##\nare structurally novel antiepileptic compounds that are highly effective anticonvulsants in animal tests (Maryanoff, B. E, Nortey, S. O., Gardocki, J. F., Shank, R. P. and Dodgson, S. P. J Med. Chem. 30, 880-887, 1987; Maryanoff, B. E., Costanzo, M. J., Shank, R. P., Schupsky, J. J., Ortegon, M. E., and Vaught J. L. Bioorganic & Medicinal Chemistry Letters 3, 2653-2656, 1993). These compounds are covered by U.S. Pat. No.4,513,006. One of these compounds 2,3:4,5-bis-O-(1-methylethylidene)-.beta.-D-fructopyranose sulfamate known as topiramate has been demonstrated in clinical trials of human epilepsy to be effective as adjunctive therapy or as monotherapy in treating simple and complex partial seizures and secondarily generalized seizures (E. FAUGHT, B. J. WILDER, R. E. RAMSEY, R. A. REIFE, L D. KRAMER, G. W. PLEDGER, R. M. KARIM et. al., Epilepsia 36 (S4) 33, 1995; S. K. SACHDEO, R. C. SACHDEO, R. A. REIFE, P. LIM and G. PLEDGER, Epilepsia 36 (S4) 33, 1995), and is currently marketed for the treatment of simple and complex partial seizure epilepsy with or without secondary generalized seizures in approximately twenty countries including the United States, and applications for regulatory approval are presently pending in several additional countries throughout the world.\nCompounds of Formula I were initially found to possess anticonvulsant activity in the traditional maximal electroshock seizure (MES) test in mice (SHANK, R. P., GARDOCKI, J. F., VAUGHT, J. L., DAVIS, C. B., SCHUPSKY, J. J., RAFFA, R. B., DODGSON, S. J., NORTEY, S. O., and MARYANOFF, B. E., Epilcpsia 35 450-460, 1994). Subsequent studies revealed that Compounds of Formula I were also highly effective in the MES test in rats. More recently topiramate was found to effectively block seizures in several rodent models of epilepsy (J. NAKAMURA, S. TAMURA, T. KANDA, A. ISHII, K. ISHIHARA, T. SERIKAWA, J. YAMADA, and M. SASA, Eur. J. Pharmacol. 254 83-89, 1994), and in an animal model of kindled epilepsy (A. WAUQUIER and S. ZHOU, Epilepsy Res. 24, 73-77, 1996).\nCluster headache is an excruciating painful disorder associated with considerable suffering which is characterized by severe, short duration, unilateral orbital-temporal pain, Ipsilateral autonomic dysfunction and chronobiological disturbances. (A. Kudrow, The pathogenesis of a cluster headache, Curr. Opin. Neurol. 7:278-282, 1994; Headache Classification Committee of the International Headache Society, Classification and diagnostic criteria for headache disorders, cranial neuralgias and facial pain, Cephalgia S7:1-96, 1988). Although treatment is often successful, some patients have intractable pain. Thus, a need remains for effective treatments for cluster headaches.\nStudies on five patients treated with topiramate have revealed previously unrecognized pharmacological properties which suggest that topiramate is effective in treating cluster headaches."}
{"text": "Many aerial vehicles (e.g., manned or unmanned vehicles such as airplanes, helicopters or other airships) are configured to operate in two or more flight modes. As one example, an aerial vehicle may be configured to engage in forward flight, or substantially horizontal flight, a mode in which the aerial vehicle travels from one point in space (e.g., a land-based point or, alternatively, a sea-based or air-based point) to another point by traveling over at least a portion of the Earth. In forward flight, the aerial vehicle may be maintained aloft by one or more net forces of lift that are typically induced by airflow passing over and below wings, consistent with a pressure gradient. As another example, an aerial vehicle may be configured to engage in vertical flight, a mode in which the aerial vehicle travels in a vertical or substantially vertical direction from one altitude to another altitude (e.g., upward or downward, from a first point on land, on sea or in the air to a second point in the air, or vice versa) substantially normal to the surface of the Earth, or hovers (e.g., maintains a substantially constant altitude), with an insubstantial change in horizontal or lateral position. In vertical flight, the aerial vehicle may be maintained aloft by one or more net forces of lift that are typically induced by rotating blades of a propeller or another source. As yet another example, an aerial vehicle may be configured to engage in both forward and vertical flight, a hybrid mode in which a position of the aerial vehicle changes in both horizontal and vertical directions.\nAn aerial vehicle that is configured to operate in multiple modes may utilize one or more propulsion systems and/or control surfaces (e.g., wings, rudders, ailerons, flaps or other components) at different times, depending on requirements of a given mission in which the aerial vehicle is to operate in each of such modes. For example, an aerial vehicle may utilize a first set of motors or rotors when operating in forward flight, and a second set of motors or rotors when operating in horizontal flight. Likewise, the aerial vehicle may utilize a first set of control surfaces when operating in horizontal flight, and a second set of control surfaces when operating in vertical flight. When motors, rotors, control surfaces or other components of an aerial vehicle are not being utilized for propulsion or control, such components merely act as dead weight to the aerial vehicle.\nThe use of imaging devices or other sensors on aerial vehicles is increasingly common. In particular, unmanned aerial vehicles, or UAVs, are frequently equipped with one or more imaging devices such as digital cameras; position sensors such as Global Positioning System, or GPS, sensors; radar sensors; or laser sensors, such as light detection and ranging, or LIDAR, sensors. Such sensors aid in the guided or autonomous operation of an aerial vehicle, and may be used to determine when the aerial vehicle has arrived at or passed over a given location, when the aerial vehicle is within range of one or more structures, features, objects or humans (or other animals), or for any other purpose. Outfitting an aerial vehicle with one or more of such sensors typically requires installing housings, turrets or other structures or features by which such sensors may be mounted to the aerial vehicle. Such structures or features add weight to the aerial vehicle, and may increase the amount or extent of drag encountered during flight, thereby exacting a substantial operational cost from the aerial vehicle for the use of such sensors in exchange for their many benefits."}
{"text": "Usually, in the CIS-based thin film solar cell module, layers including a metallic base electrode layer, a p type light absorbing layer, a high resistance buffer layer and an n type window layer (a transparent conductive film) are respectively laminated on the surface of a substrate (109) to form the CIS-based thin film solar cell module. A filler (103) having a sealing effect such as an EVA (Ethylene-Vinyl Acetate) resin, PVB (Polyvinyl Butyral), etc. is put thereon and cover glass (102) of an upper surface is laminated and attached thereon. The obtained solar cell module is surrounded by a frame (101) made of aluminum etc. to cover an end part of the solar cell module. Between the frame and the solar cell module, a resin is sandwiched (not shown in the drawing) to prevent moisture such as water from entering from the end part of the cover glass (102) and improve a weather resistance (see FIG. 1).\nOn the other hand, there is a frameless solar cell module to which an aluminum frame is not attached in order to lighten the solar cell module and reduce a production cost. As such a frameless solar cell module, a solar cell module that includes a light receiving surface side film, a light receiving surface side filler, a plurality of solar cell elements electrically connected together by connecting tabs, a back surface side filler and a back surface side film which are sequentially arranged so as to be piled, and has a structure in which a peripheral edge part of the light receiving surface side film is fusion welded to a peripheral edge part of the back surface side film is proposed (see patent document 1).\nFurther, as another frameless solar cell module, a structure in which when a frameless solar cell module is laid on a member to be attached such as a roof of a residence having a gradient, rod shaped joint fillers are held between the solar cell modules adjacent in the direction of the gradient of the member to be attached to lay the solar cell modules so that all the rod shaped joint fillers do not protrude from the surfaces of the solar cell modules is proposed (see patent document 2).\nFurther, a frameless solar cell module that has an edge space (a space where device layers are not piled) provided in a periphery of a solar cell circuit is also proposed (see FIG. 2 and patent document 3). When the edge space is provided, the frame does not need to be attached and a production cost can be more reduced and the solar cell module can be more lightened than the solar cell module of a type having the frame. As a manufacturing method of the solar cell module of this type, after a laminated film (a first electrode (108)/a semiconductor layer (107)/a second electrode (104)) is formed on the entire surface of a light receiving surface side of a substrate (109), the laminated film of an area corresponding to the edge space is removed by a laser or a sandblaster etc. to form the edge space (see FIG. 2). For instance, patent document 4 discloses a technique for removing a laminated film of an edge space area by using a YAG laser."}
{"text": "Recently linear low density polyethylenes have become available in pelletized form and they have properties which render them desirable for use in making film by the bubble extrusion process. Unfortunately high clarity, high tack film which is essential to food packaging and pallet wrap films can only be accomplished with such LLDPE resins by chill roll casting or similar techniques which employ chilled rolls or mandrels for rapid cooling of the molten extrudate."}
{"text": "A muffler usually comprises two housing parts, namely two half-shells, which are joined with a seal being interposed therebetween. Such a seal must withstand high temperatures and compensate for unevenness at the sealing faces as well as withstand vibrations and gas pulsations of the muffler.\nExperience has shown that the silicon seals or fiber seals normally used do not withstand the high temperatures which can be approximately 500.degree. C. so that considerable leakage can develop even after a short operating time. The arrangement of copper seals is likewise not satisfactory since warpage occurring because of high temperatures can lead to leakage and mechanical wear occurs because of unavoidable relative movements between the housing parts produced by vibration."}
{"text": "This invention relates generally to layered perovskite superconductors and, more particularly, to substrates on which layered perovskite superconductors are deposited to form high-frequency electronic devices and circuits.\nMaterials exhibiting superconductivity at temperatures above the advantageous liquid-nitrogen temperature of 77 K were discovered only recently and have triggered a world-wide explosion in scientific and technological research. The first material to exhibit superconductivity above the temperature of liquid nitrogen was an oxygen-depleted layered perovskite compound of yttrium, barium, copper and oxygen, identified by the formula Y1Ba2Cu3O7. Since this discovery, other similar layered perovskite copper oxide compounds, identified by the formula R1Ba2Cu3O7, where R is a rare earth element, have also been found to be superconductive at temperatures above the liquid-nitrogen temperature. This particular group of layered perovskite superconductors is commonly referred to as xe2x80x9c1-2-3xe2x80x9d compounds, because of the number of atoms of each metal element in the copper oxide compound.\nStill other layered perovskite copper oxide compounds, with even higher critical temperatures (the temperature at which superconductivity occurs), have been more recently discovered. These newer compounds contain four metallic elements instead of the three metallic elements contained in the xe2x80x9c1-2-3xe2x80x9d compounds, and they do not contain a rare earth element. In place of the rare earth element, these newer compounds contain metals such as bismuth or thallium.\nThe major advantage of the layered perovskite superconductors is that superconductive temperatures can be maintained using liquid nitrogen, which is considerably less expensive and troublesome than using liquid helium, as required in the past. Therefore, these superconductors can be expected to find many new applications. One major application already being investigated is integrated circuits, in which thin films of these new superconductors are deposited on substrates to form, for example, Josephson junctions, waveguides and microwave transmission lines. These superconductor circuit elements can be combined to form high-speed, high-frequency and low-power integrated circuits with unparalleled performance.\nHowever, thin films of the layered perovskite superconductors can only be grown with optimal properties on substrates having crystal structures and lattice constants that closely match those of the superconductors. Strontium titanate (SrTiO3) is one such material, and is currently being used as a substrate. Unfortunately, strontium titanate is unsuitable at high frequencies because it is very lossy and has an extremely high dielectric constant at superconductive temperatures. Accordingly, there has been a need for a substrate material having good high frequency characteristics and a crystal structure and lattice constant that closely match the crystal structures and lattice constants of the layered perovskite superconductors. The present invention clearly fulfills this need.\nThe present invention resides in a substrate of lanthanum aluminate (LaAlO3) on which thin films of layered perovskite copper oxide superconductors are formed. Lanthanum aluminate, with a pseudo-cubic perovskite crystal structure, has a crystal structure and lattice constant that closely match the crystal structures and lattice constants of the layered perovskite superconductors. Therefore, it promotes epitaxial film growth of the superconductors, with the crystals being oriented in the proper direction for good superconductive electrical properties, such as a high critical current density. In addition, LaAlO3 has good high frequency properties, such as a low loss tangent and low dielectric constant at superconductive temperatures. Finally, lanthanum aluminate does not significantly interact with the superconductors. Lanthanum aluminate can also be used to form thin insulating films between the superconductor layers, which allows for the fabrication of a wide variety of superconductor circuit elements.\nIt will be appreciated from the foregoing that the present invention represents a significant advance in the field of superconductors. Other features and advantages of the present invention will become apparent from the following more detailed description, taken in conjunction with the accompanying drawings, which illustrate, by way of example, the principles of the invention."}
{"text": "Various mechanisms exist for power management in a platform. However, existing power management techniques require a process or agent to have access to all components to be managed.\nExisting multi-processor platforms may have processors, devices and/or memory partitioned. One platform may then behave as two computing devices. A platform may have multiple processors with hardware partitioning so that the platform may act as multiple machines with multiple operating systems. There may be only a single motherboard on the platform.\nPower utilization of platforms is becoming more important. It is desirable to minimize or optimize the power that a platform or set of platforms utilize. As partitioning becomes more prevalent, it may become desirable to manage the power utilization of the various partitions of a system and the system overall.\nHowever, partitioning has traditionally been performed in enterprise-class machines, such as servers, and other large-size deployed machines. Power utilization has become more important as more machines are used in one location. Environmental concerns are now an issue. For instance, it may not be possible to co-locate enough compute power and maintain an appropriately cooled environment. The air conditioning unit may not be able to handle the heat dissipation of that much power utilization. With the multi-core and mini-core deployments planned for the future, many usage models require partitioning. It is becoming more common for platforms to comprise more than one processor. Further, other devices are often present in multiples within the platform.\nPortions of a system may be sequestered, for partitioning. Power management in an existing system may use an agent that can monitor the processor burden rate (busy time), control throttling, etc. The agent is typically in control of the platform, but when resources are sequestered, the agent no longer maintains control of the sequestered resources. If the agent cannot see a resource, the agent can no longer manage that resource. Thus, the traditional power management schemes are no longer viable solutions for partitioned platforms with sequestering."}
{"text": "Mission requirements for near-future radars dictate high levels of operational capability provided by systems that are light in weight. Such radars must feature agile, reconfigurable beams coupled with high effective transmit power and high receive sensitivity.\nThe operational requirements are fulfilled by adopting large aperture active phased array antennas having transmit/receive (TIR) electronics distributed with the radiating elements. Distributing the active TIR circuits over the array antenna also necessitates distributing their associated prime power converters and controllers, plus providing means for effective thermal management and conveying RF/power signals. It is desirable that these phased array antennas be realized with minimum weight to promote high mobility in ground radar applications and to minimize top-side mass for shipboard systems.\nAccordingly, there is a need for a lightweight active phased array antenna having distributed transmit/receive (T/R) electronics radiating elements, power converters, and controllers. Such a phased array antenna should also have effective thermal management and a mechanism for conveying the RF/power signals."}
{"text": "Given a data matrix X of size n by p, clustering assigns the observations (rows of X) to clusters, or groups based on some or all of the data variables (columns of X). Clustering is a cornerstone of business intelligence, with wide-ranging applications such as market segmentation and fraud detection."}
{"text": "Distributed database systems include a plurality of storage devices spread among a network of interconnected computers. The distributed database systems typically have greater reliability and availability than parallel database systems, among other benefits. Various internet services, for example social networking services, employ distributed database systems to manage the storage and retrieval of information. Generally, the need to efficiently and accurately read and write data across the database system increases with a greater amount of information, a greater number of users, and stricter latency requirements.\nIn distributed systems, there is a need for coordination and synchronization among the storage devices. In particular, the coordination and synchronization may be needed to ensure that only certain processes should perform certain tasks, ensure the proper execution of a set of functions, ensure the allocation of unique IDs, orchestrate multiple processes, and/or facilitate other coordination operations. Apache ZooKeeper™ is an example of an existing centralized service for maintaining configuration information, providing distributed synchronization, and providing group services.\nHowever, the application programing interface (API) for ZooKeeper™ is complex such that faulty function calls or faulty services interfacing with ZooKeeper™ only fail in certain edge cases. As a result, these faulty function calls or faulty services may be difficult to detect or may otherwise go undetected, may be difficult for programmers to debug, and/or may cause difficulties in fencing. Accordingly, there is an opportunity for configurations and frameworks that decrease the complexity exposed to developers for many use-cases of a distributed database."}
{"text": "The present invention generally relates to the preparation of low hysteresis rubber.\nThe hysteresis of an elastomer refers to the difference between the energy applied to deform the elastomer and the energy recovered as the elastomer returns to its initial, undeformed state. In many industries, it is often desirable to produce elastomeric compounds exhibiting reduced hysteresis when properly compounded with other ingredients, such as reinforcing agents, and then vulcanized. Such elastomers, when compounded, fabricated, and vulcanized into components for constructing articles such as tires, power belts, and the like, will manifest properties of increased rebound and reduced heat build-up when subjected to mechanical stress during normal use. In pneumatic tires, lowered hysteresis properties are associated with reduced rolling resistance and heat build-up during operation of the tire. As a result, lower fuel consumption is realized in vehicles using such pneumatic tires.\nCompounded elastomer systems are known in the art and are comprised of at least one elastomer, a reinforcing filler agent (such as carbon black or mineral fillers such as clay, silica, and the like), and a vulcanizing system. Notably, interaction between the elastomer molecules and the reinforcing filler agent(s) is known to affect hysteresis.\nIt has been recognized that hysteresis and other physical properties of compounded elastomer systems can be improved by ensuring good dispersion of carbon black throughout the elastomer component. In particular, it has been recognized that individual carbon black aggregates should be well dispersed in order to minimize the contact between them. By preparing elastomers with end groups capable of interacting with reinforcing fillers such as carbon black, dispersion of individual carbon black aggregates is stabilized, reducing interaggregate contacts and thereby reducing hysteresis. Thus, increased interaction between elastomer chain ends and reinforcing fillers will tend to decrease hysteretic losses.\nThus, there exists a need in the art for increasing interaction between elastomer molecules and reinforcing fillers within a compounded elastomer system. More particularly, there exists a need in the art for a terminator that provides elastomer molecules with end groups that interact with the functionality of carbon black and other reinforcing fillers. There also exists a need in the art for a process for terminating polymers with such desired end groups.\nIn the present invention, lithium-initiated polymers are terminated with oxazoline compounds. Useful oxazoline compounds may be selected from organic halide oxazoline compounds and branched or unbranched vinyl oxazoline compounds.\nWhen organic halide oxazoline compounds are employed, the carbon-lithium chain end of the propagating polymer reacts with the halide provided by the oxazoline compound, and the polymer is terminated with the organic oxazoline.\nWhen a branched or unbranched vinyl oxazoline compound is employed, the vinyl oxazoline compound is inserted between the propagating polymer and the lithium atom, thereby effectively terminating further polymer propagation with chemical units other than vinyl oxazoline compounds or, more broadly, monomers that will polymerize with a substituted acrylate anion, such as, for example, esters of methacrylic acid, acrylic acid, cyanoacrylic acid and the like.\nThe oxazoline ring of these terminators interacts with carbon black, and the resulting functional polymers therefore exhibit reduced hysteresis in cured, compounded rubber, when compared to unterminated, yet otherwise identical, controls.\nGenerally, the present invention provides a method of forming an oxazoline functionalized polymer. The method includes the steps of forming a solution of one or more anionically polymerizable monomers in an alkane solvent, and polymerizing the monomers in the presence of a lithio-containing initiator so as to produce a living lithium polymer as is known in the art. The living lithium polymer is then terminated with an oxazoline compound according to the formula (I): \nwherein R1 to R5 are selected according to the following criteria:\nR1 is selected from the group consisting of:\nan organic halide according to the following formula (II): \nwherein each R6 may be the same or different and is selected from the group consisting of a branched or linear C1-C20 alkyl moiety, a C3-C20 cycloalkyl moiety, a C6-C20 aryl moiety, and a C7-C20 alkylaryl moiety, R7 is selected from the group consisting of a branched or linear C1-C20 alkylene moiety, a C4-C20 cycloalkylene moiety, and a C6-C20 alkylarylene moiety, and X is a halogen;\na branched or unbranched vinyl moiety according to the following formula (III): \nwherein R8 is selected from the group consisting of hydrogen, a branched or linear C1-C20 alkyl moiety, a C3-C20 cycloalkyl moiety, a C6-C20 aryl moiety, and a C7-C20 alkylaryl moiety;\na moiety according to formula (IV): \nwherein each R6 may be the same or different and is selected according to the same criteria set forth in formula (II); and\nhydrogen; and\nR2-R5 may be the same or different and are selected from the group consisting of:\nhydrogen;\nan organic halide, xe2x80x94R7xe2x80x94X, wherein R7 is as already defined, and X is a halogen; and\na branched or linear C1-C20 alkyl moiety, a C3-C20 cycloalkyl moiety, a C6-C20 aryl moiety, and a C7-C20 alkylaryl moiety, with the proviso that if R1 is selected from either an organic halide according to formula (II) or a vinyl moiety according to formula (III), R2-R5 are selected from hydrogen, a branched or linear C1-C20 alkyl moiety, a C3-C20 cycloalkyl moiety, a C6-C20 aryl moiety, and a C7-C20 alkylaryl moiety; and if R1 is selected from either hydrogen or a moiety according to formula (IV), one of R2-R5 is selected from an organic halide, xe2x80x94R7xe2x80x94X, the other R2-R5 groups are selected from the group consisting of hydrogen, a branched or linear C1-C20 alkyl moiety, a C3-C20 cycloalkyl moiety, a C6-C20 aryl moiety, and a C7-C20 alkylaryl moiety.\nThe present invention also provides compounds according to the following formula: \nwherein X is a halogen, and each Z is the same and is selected from the group consisting of hydrogen and a methyl group.\nAs mentioned, it has been found that the oxazoline ring of the terminators disclosed herein beneficially interact with carbon black to reduce hysteresis in rubbers compounded with carbon black. Thus, the present invention also provides a method for producing a tire having reduced hysteresis. This method includes terminating a lithium-initiated rubber with an oxazoline compound; and compounding the oxazoline-terminated rubber with carbon black, wherein the oxazoline compound employed to terminate the lithium-initiated rubber is selected according to formula (I) under the criteria as set forth above with respect to formula (I)."}
{"text": "File systems organize and track where data is stored in memory and where free or available space exists in memory. Distributed or clustered file systems can store thousands or even millions of files. These files and the corresponding metadata are distributed across a large number of storage devices, such as disk drives.\nCounters are used in the file system to track counts for various types of information. Many file systems use simple counters that are not distributed. This approach does not scale well on clustered file systems. Other file systems implement counters that cannot be recovered in the event of corruption without a scan of the entire file system by a file system checker (fsck). Such scans can be quite time-consuming for large file systems.\nMethods and systems are needed to improve the management of counters in clustered file systems."}
{"text": "Wireless power transfer (WPT) or wireless energy transmission is the transmission of electrical power from a power source to a receiving device without using solid wires or conductors. Generally, the term refers to a number of different power transmission technologies that use time-varying electromagnetic fields. In wireless power transfer, a transmitter device is connected to a power source, such as the mains power line, and transmits power by electromagnetic fields across an intervening space to one or more receiver devices, where it is converted back to electric power and utilized. Wireless power techniques may fall into two categories, non-radiative and radiative.\nIn near-field or non-radiative techniques, power may be transferred over short distances by magnetic fields using inductive coupling between coils of wire by electric fields using capacitive coupling between electrodes, for example. Using these techniques, it may be possible to transfer power wirelessly within two meters distance with an efficiency of 70% at a frequency of 100 kHz. Applications of this type include, but are not limited to, cell phone, tablets, laptops, electric toothbrush chargers, RFID tags, smartcards, and chargers for implantable medical devices like artificial cardiac pacemakers, and inductive powering or charging of electric vehicles like cars, trains or buses. In radiative or far-field techniques, sometimes called “power beaming”, power may be transmitted by beams of electromagnetic radiation, like microwave or laser beams. These techniques can transport energy longer distances but are typically aimed at the receiver.\nNumerous standardized technologies have been developed over the years, including, but not limited to, Alliance for Wireless Power (“A4WP”) which is based on an interface standard (“Rezence”) for wireless electrical power transfer based on the principles of magnetic resonance, where a single power transmitter unit (PTU) is capable of charging one or more power receiver units (PRUs). The interface standard supports power transfer up to 50 Watts, at distances up to five centimeters. The power transmission frequency may be 6.78 MHz, and up to eight devices can be powered from a single PTU depending on transmitter and receiver geometry and power levels. A Bluetooth Smart link may be provided in an A4WP system for controlling power levels, identification of valid loads, and protection of non-compliant devices.\nAnother standardized WPT technology includes “Qi” which relies on electromagnetic induction between planar coils. A Qi system may include a base station, which is connected to a power source and provides inductive power, and a device that consumes inductive power. The base station may include a power transmitter having a transmitting coil that generates an oscillating magnetic field; the device may include a power receiver holding a receiving coil. The magnetic field from the transmitter may induce an alternating current in the receiving coil by Faraday's law of induction. A further standardized WPT technology includes “Powermat” adopted by the Power Matters Alliance (PMA), which is based upon inductively coupled power transfer, where a transmitter may vary current in a primary induction coil to generate an alternating magnetic field from within a charging spot. The receiver in the consuming device may have its own induction coil that takes power from the magnetic field and converts it back into electrical current to charge the device battery. An additional part of the technology is the use of system control communication via Data over Coil (DoC), where the receiver may send feedback to the transmitter by changing the load seen by the transmitter coil. The protocol is frequency based signaling, thus enabling fast response of the transmitter.\nIn WPT systems, it is often necessary to place a consuming device, having a receiver antenna, within the general area and range of a WPT transmitter in order to receive power. While adequate power may be provided when the consuming device is physically placed in the general area and range of the WPT transmitter, the placement may not be optimal, which may result in less-than optimal power transfer. In cases where the receiving entity is of a larger size (e.g., car, truck, train, etc.) this problem may be compounded. Accordingly, technologies are needed to optimize wireless charging alignment between a WPT transmitter and one or more receivers."}
{"text": "A. Technical Field\nThe present invention relates to an ink jet recording material and a coating agent therefor, more particularly, relates to: an ink jet recording material which displays sufficient water resistance and further, excellent definition of initial images; and a coating agent therefor to give such an ink jet recording material.\nInk jet recording apparatuses are machines such as printers, facsimile machines, and copiers. For the ink jet recording apparatus, U.S. Pat. No. 5,486,854 issued Jan. 23, 1996 is hereby incorporated by reference. The ink jet recording apparatus is to make a recording by jetting ink from a recording means (i.e. a recording head which is a part of the ink jet recording apparatus) onto a recording material, wherein the recording material is, for example, a paper sheet, or a plastic sheet such as transparent PET sheet.\nThe ink jet recording material, generally, suffers from a lack of water resistance. That is, an exposure to water usually will dissolve and destroy the image (the imaged ink). To prevent this, the image must be rendered water-resistant, and further, if the ink is, for example, an organic dye, then the image must also be fixed. Examples of known methods of giving the image the water resistance to thereby fix the image include arts to fix dyes with mordants and arts involving the use of adsorptive pigments. However, operations thereof are complicated, or the optimal method is different according to the composition of the ink, so the above known methods are not commonly usable means.\nIn comparison, JP-A-035090/1998 discloses a water-resistant composition for ink jet recording sheets, comprising a polymer containing an amino group (and/or quaternary ammonium salt) and a carboxyl group (and/or acid anhydride) and a crosslinking agent containing at least two oxazoline groups, as a method of not giving the image the water resistance, but carrying out a water-resisting pretreatment for the recording material. However, as to this technique, not only is the resultant water-resistification insufficient, but also there are problems on the definition of initial images.\nAn object of the present invention is to provide: an ink jet recording material which displays sufficient water resistance and further, excellent definition of initial images; and a coating agent therefor to give such an ink jet recording material.\nTo solve the above problems, the present invention provides the following:\nA first coating agent for ink jet recording materials, according to the present invention, comprises: an aqueous polymer having a carboxyl group; and a water-soluble polymer having an oxazoline group as a crosslinking agent.\nA second coating agent for ink jet recording materials, according to the present invention, comprises a polymer and a crosslinking agent, wherein the polymer has both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group.\nAn ink jet recording material, according to the present invention, has on at least one face thereof a coated and cured layer of the above present invention first coating agent and/or a coated and cured layer of the above present invention second coating agent.\nThese and other objects and the advantages of the present invention will be more fully apparent from the following detailed disclosure.\nThe first coating agent, according to the present invention, comprises: an aqueous polymer having a carboxyl group; and a water-soluble polymer having an oxazoline group as a crosslinking agent.\nThe aqueous polymer having a carboxyl group is not especially limited if it is a polymer that has a carboxyl group and further is aqueous (that is, water-soluble or water-dispersible). To obtain the polymer having a carboxyl group, a monomer having a carboxyl group is, for example, polymerized as a raw material, or a carboxyl group is introduced into a polymer (as prepared beforehand) by denaturation. To obtain the aqueous polymer, a hydrophilic monomer (available whether it has a carboxyl group or not) is, for example, used in the ratio of preferably 10 mol % or higher, more preferably 50 mol % or higher, to the entire monomer component.\nExamples of the monomer having a carboxyl group include: unsaturated monocarboxylic acids, such as acrylic acid, methacrylic acid, and crotonic acid; unsaturated dicarboxylic acids, such as maleic acid, itaconic acid, citraconic acid, and fumaric acid; and unsaturated dicarboxylic anhydrides, such as maleic anhydride, itaconic anhydride, and citraconic anhydride.\nExamples of the method, in which a carboxyl group is introduced into a polymer (as prepared beforehand) by denaturation, include a method including the step of jumping up a polymer, having an OH group in opposite terminal portions, with pyromellitic dianhydride.\nExamples of the hydrophilic monomer include: monomers having a carboxyl group; and other monomers, such as hydroxyethyl acrylate, hydroxyethyl methacrylate, vinylpyrrolidone, dimethylaminoethyl acrylate, and chloridized-triaminoethyl methacrylate.\nExamples of the aqueous polymer having a carboxyl group, as preferably usable in the present invention, include: polyvinyl alcohols having a carboxyl group (for example, anionic KEPS series made by Dai-ichi Kogyo Seiyaku Co., Ltd., and K Polymer made by Kuraray Co., Ltd.); (meth)acrylic ester copolymers (for example, Arolon made by Nippon Shokubai Co., Ltd.); vinyl ether-maleic anhydride copolymers (Gantrez AN series made by ISP); and aqueous polymers having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group. Particularly preferable ones are the aqueous polymers having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group.\nExamples of the above aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group include a polymer having a structural unit of general formula (1) below: \nwherein: each of R1 and R2, independently of each other, denotes a\nhydrogen atom or a methyl group;\nR3 denotes a divalent organic residue;\nR4 denotes a monovalent organic residue;\nZ denotes a halogen atom;\na denotes an integer of 1xcx9c1,000;\nb denotes an integer of 3xcx9c110,000; and\nn denotes an integer of 3xcx9c5,000.\nThis polymer can be synthesized by carrying out cationic polymerization of an oxazoline compound in the presence of an unsaturated halide to synthesize a polyoxazoline macromonomer having a radical-polymerizable double bond at a polymerization-initiating terminal, and then copolymerizing this polyoxazoline macromonomer and a monomer having a carboxyl group, as is illustrated by the following chemical reaction formula: \nwherein R1, R2, R3, R4, Z, a, b, and n are the same as those in general formula (1).\nExamples of the oxazoline compound, as usable in the above cationic polymerization, include 2-methyl-2-oxazoline, 2-ethyl-2-oxazoline, 2-(n-propyl)-2-oxazoline, 2-(i-propyl)-2-oxazoline, 2-(n-butyl)-2-oxazoline, 2-(i-butyl)-2-oxazoline, and 2-(t-butyl)-2-oxazoline. Preferable ones among them are compounds with not more than 3 carbon atoms in R4, of which specific examples include 2-ethyl-2-oxazoline.\nExamples of the above unsaturated halide include chloromethylstyrene, allyl chloride, epichlorohydrin, and chloroethyl vinyl ether. A preferable one among them is chloromethylstyrene.\nAs to the monomer having a carboxyl group, compounds in which R2 is a hydrogen atom or a methyl group are preferable among the above-exemplified monomers having a carboxyl group. Specific examples of such compounds include acrylic acid and methacrylic acid.\nTherefore, preferable examples of the aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group include poly(2-methyl-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-ethyl-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(n-propyl)-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(i-propyl)-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(n-butyl)-2-oxazoline)/(meth)acrylic acid copolymers, poly(2-(i-butyl)-2-oxazoline)/(meth)acrylic acid copolymers, and poly(2-(t-butyl)-2-oxazoline)/(meth)acrylic acid copolymers.\nThe above aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group is available whether it contains a structural unit other than the structural unit of general formula (1) above or not.\nIn addition, in the above general formula (1), a is an integer of 1xcx9c1,000, preferably 1xcx9c500, and b is an integer of 3xcx9c10,000, preferably 3xcx9c5,000, and n is an integer of 3xcx9c5,000, preferably 10xcx9c5,000, more preferably 2xcx9c100, still more preferably 5xcx9c500.\nThe weight-average molecular weight of the above aqueous polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group is preferably in the range of 50,000xcx9c1,000,000, more preferably 100,000xcx9c500,000.\nExamples of the water-soluble polymer having an oxazoline group as a crosslinking agent, as contained in the first coating agent according to the present invention, include a polymer containing an oxazoline group as obtained by polymerizing a monomer component which comprises an addition-polymerizable oxazoline and, if necessary, further comprises a monomer copolymerizable therewith. Examples of the addition-polymerizable oxazoline include 2-vinyl-2-oxazoline, 2-vinyl-4-methyl-2-oxazoline, 2-vinyl-5-methyl-2-oxazoline, 2-isopropenyl-2-oxazoline, 2-isopropenyl-4-methyl-2-oxazoline, 2-isopropenyl-4-ethyl-2-oxazoline, 2-isopropenyl-5-methyl-2-oxazoline, 2-isopropenyl-5-ethyl-2-oxazoline, and 2-isopropenyl-4,5-dimethyl-2-oxazoline. These may be used either alone respectively or in combinations with each other. A preferable one among them is 2-isopropenyl-2-oxazoline, because it is industrially easily available.\nThe amount of the addition-polymerizable oxazoline, as used, is not especially limited, but is preferably 5 weight % or larger, more preferably in the range of 30xcx9c60 weight %, of the entire monomer component. In the case where the amount is smaller than 5 weight %, the extent of the curing is insufficient. In the case where the amount exceeds 60 weight %, it will have a bad effect on the water resistance.\nExamples of the monomer copolymerizable with the addition-polymerizable oxazoline include: (meth)acrylic esters, such as methyl (meth)acrylate, butyl (meth)acrylate, 2-ethylhexyl (meth)acrylate, methoxypolyethylene glycol (meth)acrylate, polyethylene glycol mono(meth)acrylate, 2-hydroxyethyl (meth)acrylate, and 2-aminoethyl (meth)acrylate and its salts; unsaturated nitriles, such as (meth)acrylonitrile; unsaturated amides, such as (meth)acrylamide, N-ethylol(meth)acrylamide, and N-(2-hydroxyethyl)(meth)acrylamide; vinyl esters, such as vinyl acetate and vinyl propionate; vinyl ethers, such as methyl vinyl ether and ethyl vinyl ether; xcex1-olefins, such as ethylene and propylene; halogen-containing xcex1,xcex2-unsaturated monomers, such as vinyl chloride, vinylidene chloride, and vinyl fluoride; and xcex1,xcex2-unsaturated aromatic monomers, such as styrene, xcex1-methylstyrene, and sodium styrenesulfonate. These may be used either alone respectively or in combinations with each other.\nTo obtain the water-soluble polymer, the ratio of the hydrophilic monomer to the entire monomer component to be polymerized is preferably 50 weight % or higher, particularly preferably in the range of 60xcx9c90 weight %. Examples of the hydrophilic monomer include addition-polymerizable oxazolines, 2-hydroxyethyl (meth)acrylate, methoxypolyethylene glycol (meth)acrylate, polyethylene glycol mono(meth)acrylate, and 2-aminoethyl (meth)acrylate and its salts, and further, sodium (meth)acrylate, ammonium (meth)acrylate, (meth)acrylonitrile, (meth)acrylamide, N-methylol(meth)acrylamide, N-(2-hydroxyethyl)(meth)acrylamide, and sodium styrenesulfonate, as are selected from among the aforementioned monomer components.\nThe aqueous polymer having a carboxyl group is contained in the first coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 50xcx9c90 weight %, in terms of solid content, to the entire weight of the first coating agent. The water-soluble polymer having an oxazoline group, as used as the crosslinking agent, is contained in the first coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 10xcx9c50 weight %, in terms of solid content, to the entire weight of the first coating agent.\nThe first coating agent may further comprise components other than the aqueous polymer having a carboxyl group and the water-soluble polymer having an oxazoline group as the crosslinking agent, if necessary. Examples of such other components include: curing catalysts, such as paratoluenesulfonic acid (PTSA); organic or inorganic fine particles; dye mordants; pigments; dispersants; and ultraviolet absorbing agents.\nThe second coating agent, according to the present invention, comprises a polymer and a crosslinking agent, wherein the polymer has both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group. Incidentally, this xe2x80x9cpolymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl groupxe2x80x9d is mentioned in detail above.\nExamples of the crosslinking agent, as used for the second coating agent, include the water-soluble polymer having an oxazoline group as mentioned in detail above, and further, melamine, aziridine, isocyanate, and epoxy.\nThe polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group is contained in the second coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 50xcx9c90 weight %, in terms of solid content, to the entire weight of the second coating agent. The crosslinking agent is contained in the second coating agent in the ratio of preferably 5xcx9c95 weight %, more preferably 10xcx9c90 weight %, still more preferably 10xcx9c50 weight %, in terms of solid content, to the entire weight of the second coating agent.\nThe second coating agent may further comprise components other than the polymer having both a structural unit, as formed by ring-opening polymerization of an oxazoline compound, and a carboxyl group and the crosslinking agent, if necessary. Examples of such other components include: curing catalysts, such as paratoluenesulfonic acid (PTSA); organic or inorganic fine particles; dye mordants; pigments; dispersants; and ultraviolet absorbing agents.\nThe ink jet recording material, according to the present invention, has on at least one face thereof a coated and cured layer of the present invention first coating agent and/or a coated and cured layer of the present invention second coating agent.\nThe ink jet recording material is a sheet-shaped material as used to record images thereon with ink jet recording apparatuses. Examples of such ink jet recording materials include: paper, synthetic paper such as Tyvek (made by E.I. Du Pont DE NEMOURS and Co., Ltd.); cloths, such as canvas, clothing fabrics and non-woven composites; films or sheets of plastics such as polyvinyl chloride, polypropylene, and polyethylene terephthalate (PET).\nThe amount of the present invention first or second coating agent, as coated, is preferably in the range of 3xcx9c50 g, more preferably 5xcx9c40 g, per square meter. In the case where the amount is smaller than 3 g, water cannot sufficiently be absorbed from ink. In the case where the amount exceeds 50 g, much time and energy are necessary for drying the sheet. In addition, the coating thickness is preferably in the range of 1xcx9c50 xcexcm, more preferably 5xcx9c40 xcexcm. In the case where the coating thickness is less than 1 xcexcm, water cannot sufficiently be absorbed from ink. In the case where the coating thickness exceeds 50 xcexcm, the improvement of the ink absorbency cannot be expected very much, so there are economical disadvantages.\nThe curing temperature is preferably in the range of 50xcx9c200xc2x0 C., more preferably 80xcx9c150xc2x0 C. The curing time depends on the curing temperature, but is preferably in the range of 1xcx9c60 minutes, more preferably 1xcx9c30 minutes.\n(Effects and Advantages of the Invention):\nCoating an ink jet recording material with the present invention coating agent for ink jet recording materials can give an ink jet recording material which displays sufficient water resistance and further, excellent definition of initial images.\nHereinafter, the present invention is more specifically illustrated by the following examples of some preferred embodiments in comparison with comparative examples not according to the invention. However, the invention is not limited to the below-mentioned examples. In addition, in the examples, unless otherwise noted, the units xe2x80x9c%xe2x80x9d and xe2x80x9cpart(s)xe2x80x9d denote those by weight.\n less than Synthesis of FX-AA as an Aqueous Polymer having a Carboxyl Group and further a Structural Unit of General Formula (1) greater than \nA mixture of 297 g of 2-ethyl-2-oxazoline, 9.2 g of chloromethylstyrene (mixture of m- and p-isomers) and 252.4 g of ethanol was charged into a 1-liter autoclave, and then heated at 130xc2x0 C. for 4 hours, thus obtaining an ethanol solution of poly(2-ethyl-2-oxazoline) macromonomer with a styrene functional group at a polymerization-initiating terminal. Then, 43.2 g of acrylic acid, 3.5 g of 2,2xe2x80x2-azobis(isobutyronitrile) and 50 g of ethanol were added to this ethanol solution, and the resultant mixture was heated within the range of 105 to 135xc2x0 C. for 6 hours. The reaction mixture was cooled, thus obtaining an ethanol solution of a poly(2-ethyl-2-oxazoline) macromonomer/acrylic acid copolymer (655.3 g of FX-AA; solid content=59.6 weight %).\n less than Synthesis of Water-soluble Polymer (A) having an Oxazoline Group as a Crosslinking Agent greater than \nA mixture of 29 parts of deionized water and 1 part of V-50 (polymerization initiator made by Wako Pure Chemical Industries, Ltd.: 2,2xe2x80x2-azobis(2-amidinopropane) dihydrochloride) were charged into a flask having a stirrer, a reflux condenser, a nitrogen-introducing tube, a dropping funnel, and a thermometer, and then heated to 60xc2x0 C. under a slow nitrogen gas current. Thereto, a monomer mixture was added from the dropping funnel over a period of 1 hour, with this monomer mixture having been prepared beforehand and comprising 0.4 parts of ethyl acrylate, 5.6 parts of methyl methacrylate, 4 parts of methoxypolyethylene glycol methacrylate (NK Ester AM-90G made by Shin-Nakamura Chemical Industrial Co., Ltd.), and 10 parts of 2-isopropenyl-2-oxazoline. During the reaction, nitrogen gas was continuously run, and the temperature in the flask was kept at 60xc2x11xc2x0 C. After the end of the addition, the same temperature was maintained for 9 hours, and then the mixture was cooled, thus obtaining an aqueous solution of a polymer (water-soluble polymer (A)) containing a 2-oxazoline group, of which the nonvolatile content was 40%, and the pH was 8.3."}
{"text": "Recent rapid advances in the development of data processors and software implementations are making more and more data processing services available to users. Basic services, such as the processing and manipulation of data in response to instructions and commands and the managing of data including storing such data in records and files, modifying the stored data and reading the stored data are now within economic reach of both small and large users. Typically, the processing is provided by one or more digital computers cooperating with externally connected user devices, such as data terminals, which may provide the data to be manipulated, managed or transferred or the commands and instructions for the data.\nSome \"users\" of data processing services may find it preferable to process data in their own, on-premises, host computer. A single such \"user\" may have widely geographically separate \"subsidiary\" locations (such as \"outlets\" for a manufacturer) and thus have need for processing services at these various locations. One solution is to send all raw data through the common carrier network to the \"user\" host and return processed data, as needed, to terminals at the \"outlet locations\". Another solution is to provide distributed processing services at or near the various terminal locations to reduce the carrier network communications between each terminal and the host computer. More specifically, in accordance with the latter solution, it is advantageous to provide local services such as editing of \"text\", interpretation of commands, creation of forms, accumulation and manipulation of numerical and text data and the like, at locations adjacent or near to the several terminals. The carrier network communication between the host and the terminals is then limited to data after such accumulation, manipulation and editing at the various locations. It has been suggested that the provision of such distributed processing services can be appropriately handled by a service vendor for a plurality user \"customers\", reducing the cost for each \"customer\" by time-sharing the usage of processing equipment among such plurality of \"customers\".\nSince each customer has unique and differing processing needs, individualized programs, specially tailored for these needs, must be available for the several customers. It is therefore clear that each distributed processing location (or node) for the above-described service necessarily must accommodate a large number of programs (for a plurality of customers) and each node might preferably comprise a plurality of processors or CPUs, each processor supporting a plurality of programs. As is well known in the art, only one program can actually be \"executing\" at any given time in a computer. However, if the program currently running should terminate, time-out, or go into a \"wait\" state pending the occurrence of some \"event\", such as the completion of any input/output operation, another program is immediately available to begin executing, thus permitting different customer programs the opportunity to time-share computer time.\nIn discussing the execution of a program, it is important to understand the concept of \"processing\". A process is the total environment in which a program (or image) executes and has, in addition to the image, three further components, namely: (1) memory space allocated to the process for storing the image, for storing data generated and used by the image and for storing control information; (2) a software context which includes such pieces of information such as priorities and privileges allocated to the process, a unique process identifier, accounting information and so on; and (3) a hardware context which comprises the contents at any instant of time of the general purpose registers within the processor or CPU. Whenever execution of the process is suspended, its hardware context is stored in memory before another process begins execution. When a process is again allowed to execute, its hardware context is restored in the registers so that the process can resume execution at the point it was interrupted without again processing previously manipulated data or reexecuting commands or instructions.\nAs previously noted, each customer of a service vendor typically has needs and requirements that differ from, in part, and are similar to, in part, the needs and requirements of other customers. To satisfy these diverse customers, it is advantageous for the service vendor to create an individual process to execute each \"customized\" program developed for the unique needs of each customer and create other processes to execute programs having \"universal capabilities\" which satisfy common needs of a plurality of customers.\nA communication subsystem for transferring data between the several processes is a critical element in an integrated distributed-processing system suitable for diverse time-sharing customers or users. An important attribute in such a communication subsystem is to identify and select process-to-process connections in response to a request from an originating process to call a terminating end-point process which may be in a different processor in the same or in a different node and to provide communication criteria or parameters for the call, such as a rate of data flow and a grade of service. It is therefore advantageous to pass data defining such communication parameters over external communication links between the processors and the nodes (by way of foreign systems such as common carrier networks) so that the end point (users) can process such communication parameter data and intercommunicate within the parameters.\nIn order to communicate over another (foreign) system, an interface is required to convert the \"processed\" data to information suitable for conveyance over the \"foreign\" system. This interface which provides the suitable conversion is necessarily transparent to the communication parameter data provided by (and destined to) the end-point users. This limits the usefulness of the interface since it must be \"blind\" to the parameter data and thus cannot receive calls as would an end-point user or process. It is an object of this invention to enhance the usefulness of an interface process."}
{"text": "1. Field\nThis disclosure generally relates to repeaters in wireless communication systems.\n2. Background\nWireless communication systems and techniques have become an important part of the way we communicate. However, providing coverage can be a significant challenge to wireless service providers. One way to extend coverage is to deploy repeaters.\nIn general, a repeater is a device that receives a signal, amplifies the signal, and transmits the amplified signal. FIG. 1 shows a basic diagram of a repeater 110, in the context of a cellular telephone system. Repeater 110 includes a donor antenna 115 as an example network interface to network infrastructure such as a base station 125. Repeater 110 also includes a server antenna 120 (also referred to as a “coverage antenna”) as a mobile interface to mobile device 130. In operation, donor antenna 115 is in communication with base station 125, while server antenna 120 is in communication with mobile devices 130.\nIn repeater 110, signals from base station 125 are amplified using forward link circuitry 135, while signals from mobile device 130 are amplified using reverse link circuitry 140. Many configurations may be used for forward link circuitry 135 and reverse link circuitry 140.\nThere are many types of repeaters. In some repeaters, both the network and mobile interfaces are wireless; while in others, a wired network interface is used. Some repeaters receive signals with a first carrier frequency and transmit amplified signals with a second different carrier frequency, while others receive and transmit signals using the same carrier frequency. For “same frequency” repeaters, one particular challenge is managing the feedback that occurs since some of the transmitted signal can leak back to the receive circuitry and be amplified and transmitted again.\nExisting repeaters manage feedback using a number of techniques; for example, the repeater is configured to provide physical isolation between the two antennae, filters are used, or other techniques may be employed."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory device, and particularly to a semiconductor memory device suitable for mounting on a portable terminal.\n2. Description of the Background Art\nSemiconductor memory devices, which are used in portable terminals such as cellular phones, employ pseudo-SRAMs (Static Random Access Memories) for achieving large capacities and easy control. The pseudo-SRAM includes DRAM (Dynamic Random Access Memory) cells as internal memory cells, and also includes an interface, which is asynchronous to a clock similar to that of a SRAM, as an external interface for defining input control signals and address signals. A refresh operation is not controlled by an externally applied signal, but is performed internally and automatically (e.g., see Japanese Patent Laying-Open No. 2002-352577).\nFor further achieving a high speed operation, a synchronous pseudo-SRAM, which additionally employs a synchronous interface, has been available, e.g., as CellularRAM® disclosed on the Internet Web site (URL: http://www.micron.com/products/psram/cellularram/). This synchronous pseudo-SRAM includes a synchronous interface synchronized with a clock similar to that of a SRAM in addition to an asynchronous interface not synchronized with the clock.\nAccording to the synchronous pseudo-SRAM, a WAIT signal is issued to notify that external access is prohibited, before elapsing of a command latency CL, which is the number of clocks issued until data is output after reception of a read or write request, and while self-refresh is being performed. In a conventional structure, command latency CL is externally applied. However, the synchronous pseudo-SRAM does not require such external application because the WAIT signal can be internally issued to notify that the access is impossible. Command latency CL of a more appropriate value can be determined by internally determining it within the synchronous pseudo-SRAM according to a current state, as compared with the case of externally determining the command latency.\nIn contrast to the above, Japanese Patent Laying-Open No. 2001-155484 has disclosed a latency determining circuit, which can adjust a latency according to a clock frequency. This latency determining circuit includes a latency determining instruction input unit, which issues an internal start signal in response to activation of a latency determination start signal for starting latency determination in synchronization with a clock signal, a latency section defining circuit, which issues a predetermined latency determination section signal at every edge of the clock signal in response to the latency determination start signal, a delay unit issuing a delay signal to adjust the latency by delaying the internal start signal by a predetermined time, and a latency instructing circuit, which determines a latency number in response to the latency determination start signal and the latency determination section signal at the point in time of activation of the delay signal. The latency determining circuit described above can internally determine the latency without an external instruction.\nFor example, Japanese Patent Laying-Open No. 07-226077 has disclosed a method, in which self-refresh is performed simultaneously with reading or writing in the case where a pseudo-SRAM has multiple banks. More specifically, according to the paragraph [0020] in Japanese Patent Laying-Open No. 07-226077, even when one of first to fourth cell banks 10-13 is in a refresh mode, first to fourth row address latch circuits and buffers 80-83 can transmit second load address signals A0-An-2 latched by a row address latch circuit 70 to remaining three cell banks to allow access to data even during the refresh mode.\nHowever, the following problems arise in the methods disclosed in the foregoing references.\nAlthough Japanese Patent Laying-Open No. 2001-155484 has disclosed a manner of internally setting the latency, it has not disclosed timing for setting the latency.\nAlthough Japanese Patent Laying-Open No. 07-226077 has disclosed a manner of performing the refresh simultaneously with reading or writing when the bank to be refreshed is different from the bank of a reading or writing target, it has not disclosed a manner of determining the bank to be refreshed such that the bank to be refreshed may not overlap with the bank of the reading or writing target.\nFurther, in a structure having a plurality of pseudo-SRAMs accommodated in a single package, each pseudo-SRAM performs the refreshing according to timing different from the timing of the others. This results in a conflict of the WAIT signals issued from the respective pseudo-SRAMs.\nFurther, there is such a problem that data cannot be output according to timing earlier than the timing defined by command latency CL.\nIn the case where the pseudo-SRAM has multiple modes such as a synchronous mode and an asynchronous mode, and operates in one mode selected from these modes, a preamplifier, which can operate appropriately in a certain mode, may not operate appropriately in another mode, and thus may not correctly perform amplification.\nIn the case where a processing target changes to a next row during continuous reading or writing, a byte mask signal must be applied after an externally applied WAIT signal notifies of the completion of processing for starting processing of a next row (i.e., deactivation of a word line in the last processed row, activation of a word line in the next target row, amplification by a sense amplifier and others).\nFor precharging a bit line pair, a chip enable signal must be externally deactivated, and thus external control is required.\nSemiconductor memory devices such as a CellularRAM® cannot operate synchronously with a clock of a low frequency.\nWhen a chip is on standby, or is inactive because another chip is being accessed, an input/output buffer continues an operation, and thus wastes a current."}
{"text": "A combustor of a gas turbine combustion engine often includes several individual combustor cans. Within each can there are multiple swirlers which impart rotational movement to the air-fuel mixture flowing through it. A conventional configuration includes eight main swirlers and a central pilot swirler, where all swirlers have parallel axes. Compressed air flows, into each main swirler individually and into the central pilot swirler individually. Fuel is added to the air as it flows through the swirler, resulting in an air-fuel mixture flowing through each main swirler. Accordingly, in a configuration with eight main swirlers and a central pilot swirler, there are nine air-fuel mixture flows; one through each of the eight main swirlers, and one through the central pilot swirler. Each air-fuel mixture flows axially, centered on the same axis as the swirler through which it is flowing. A swirler then imparts a rotation to this axial flow, such that the air-fuel mixture exiting an individual swirler is flowing along the central axis of that swirler while simultaneously rotating around that central axis. Each of the main swirlers in this relevant configuration imparts a clockwise rotation to the air-fuel mixture flowing through it as viewed looking downstream, and the central pilot swirler imparts a counterclockwise rotation. Consequently, because each main swirler imparts a clockwise rotation to the air-fuel mixture flowing through it, the tangential velocities of the rotation of adjacent air-fuel flows will be opposite where the adjacent air-fuel flows meet. Friction in these areas where adjacent tangential fuel flows oppose each other results in shear and vortices.\nThe formation of oxides of nitrogen NOx is correlated to the temperature of combustion. Therefore, NOx emissions are reduced by reducing the temperature and size of the hot zones within the combustor. In the combustor configuration described above, the air-fuel flow through the pilot swirler runs relatively rich, i.e. a higher concentration of fuel in this mixture exists than exists in the main swirler flows. This provides a hot central flame to stabilize the overall combustor dynamics, which is necessary because the outer swirlers are unable to stabilize on their own due to the lean air-fuel mixture flowing through them. Thus, reducing NOx emissions in this configuration means reducing the size of the central pilot zone, and/or reducing the temperature of the air-fuel flow in the central pilot zone by reducing the amount of fuel in that air-fuel mixture. However, as the central pilot zone air-fuel flow (and associated temperature) is reduced, combustion dynamics (i.e. pressure oscillations) increase. These dynamic pressure oscillations can be harmful to the combustion chamber.\nDynamic pressure oscillations are associated with either the lean flammability limit of the air-fuel mixture, or fluctuations in the heat release rate of the combustion flame. Oscillations associated with the lean flammability limit are typically characterized by frequencies below 50 hertz. Oscillations associated with combustion flame heat release rate are typically associated with higher frequencies, and they and are often the limiting dynamic in the higher firing-temperature applications currently under development. High frequency pressure oscillations cause fluctuations in the heat release rate of the combustion flame, which is responsive to changes in pressure. A change in the heat release rate of the combustion flame produces pressure oscillations, and the feedback cycle repeats.\nAs a result, the ability to reduce NOx and CO emissions in the above described combustor configuration is limited by the need to minimize high frequency pressure oscillations. Accordingly, once the temperature of the central pilot zone is reduced to the point where combustion dynamics have reached a maximum safe level, NOx and CO emissions can not be reduced any more.\nConventional swirlers also have variable fuel-hole injection patterns to enable a center rich concentration of fuel in the air fuel mixture. Other patterns known in the art result in air-fuel mixtures where the fuel is either uniformly distributed throughout the air-fuel flow, or is concentrated in the outer portion of the air-fuel flow, result in high levels of combustion driven oscillations. However, because the fuel is concentrated in the center of the air-fuel flow, the peak temperature of the burn at the center of the flow is greater than the temperature of the burn of an evenly distributed air-fuel flow. This center-rich fuel configuration results in greater NOx and CO production, due to the exponential nature of NOx production with temperature.\nFurther, when main swirler flows have a center rich fuel configuration, complete combustion, which requires complete mixing of the central pilot flow and the main swirler flows, requires more time, resulting in a longer central flame, and yet further increased NOx and CO production.\nThus, there exists a need in the art to further reduce the temperature and/or size of the central pilot zone without increasing combustion dynamics, in an effort to reduce NOx and CO emissions."}
{"text": "Modern networks, such as access and aggregation networks, are moving towards Ethernet as the universal medium. However, Ethernet is not quite robust enough for carrier transport applications. Thus, a new variation termed “Carrier Ethernet” has been created. Carrier Ethernet is configured to provide Ethernet-based operations, administration, and maintenance (OAM). There are many variations of this type of Ethernet, for example Optical Transport Network (OTN) carriage of Ethernet is a form of carrier Ethernet. OTN can offer a form of tunneling, and in addition offers hardening features such as embedded forward error correction (FEC).\nReferring to FIG. 1, an Ethernet frame 10 with Virtual Local Area Network (VLAN) tagging is illustrated in one instance of Carrier Ethernet. VLAN Tagging is defined in IEEE 802.1Q as a standard to allow multiple bridged networks to transparently share the same physical network link without leakage of information between networks (i.e. trunking). The Ethernet frame 10 in IEEE 802.1Q is not actually encapsulated. Instead, the EtherType value (for Ethernet II framing) in the Ethernet header is set to hex 8100, identifying this frame as an 802.1Q frame. Also, four extra bytes are added after the Ethernet header consisting of two-byte Tag Control Information (TCI). VLAN ID (VID) is a 12-bit field specifying the VLAN to which the frame 10 belongs. A value of “0” means the frame 10 does not belong to any VLAN (so that the 802.1Q header specifies only a priority), a value of “1” is used with bridges, and a value of hex FFF is reserved for implementation use; all other values can be used as VLAN identifiers, allowing up to 4093 VLANs.\nVLAN tagging is a Media Access Control (MAC) option that provides three important capabilities not previously available to Ethernet network operators and users. First, it provides a quality-of-service (QoS) mechanism to expedite time-critical network traffic by setting transmission priorities for outgoing frames 10. Second, it allows stations to be assigned to logical groups to communicate across multiple LANs as though they were on a single LAN. For example, bridges and switches filter destination addresses and forward VLAN frames only to ports that serve the VLAN to which the traffic belongs. Finally, VLAN tagging simplifies network management and makes adds, moves, and changes easier to administer.\nThe VLAN header includes two fields: a reserved 2-byte type value, indicating that the frame is a VLAN frame, and a two-byte Tag-Control field 14 that contains both the transmission priority (0 to 7, where 7 is the highest) and a VLAN ID 12 that identifies the particular VLAN over which the frame is to be sent. The receiving MAC reads the reserved type value, which is located in the normal Length/Type field position, and interprets the received frame as a VLAN frame. If the MAC is installed in a switch port, the frame is forwarded according to its priority level to all ports that are associated with the indicated VLAN identifier. If the MAC is installed in an end station, it removes the 4-byte VLAN header and processes the frame in the same manner as a basic data frame. VLAN tagging requires that all network nodes involved with a VLAN group be equipped with the VLAN option.\nReferring to FIG. 2, an Ethernet frame 20 for Provider Backbone Transport (PBT) according to IEEE 802.1ah is illustrated in another instance of Carrier Ethernet. PBT is a set of enhancements to Ethernet technology that allows the use of Ethernet as a carrier-class transport network. PBT uses the concepts of VLAN tagging as per IEEE 802.1Q, Q-in-Q as per IEEE 802.1ad and MAC-in-MAC as per IEEE 802.1ah (Provider Backbone Bridges (PBB)) but disables the concept of flooding/broadcasting and spanning tree protocol (SPT). The idea here is to use Ethernet for connection oriented purpose as is the case with present SDH/SONET transport by stripping down the complexity involved with the present Ethernet LAN. It simplifies the OAM, as in SDH/SONET world, by using additional extensions based on IEEE 802.1ag. It also provides extensions so as to provide path protection levels similar to the UPSR protection in SDH/SONET network.\nIn the Ethernet frame 20, the tunnel is encoded by the destination MAC address of the backbone egress switch (B-DA) as well as a 12-bit VLAN-tag (backbone tag, B-VID). PBT forms a topology of B-DA rooted trees and an independent sink-tree is configured for each <B-DA, B-VID> pair. Since no SPT algorithm has to be performed, the trees need not be spanning. Thus, up to 4096 different trees can be configured for one B-DA.\nIEEE 802.1ad (Provider Bridges) is an amendment to IEEE standard IEEE 802.1Q-1998 (also known as Q-in-Q or Stacked VLANs), intended to develop an architecture and bridge protocols to provide separate instances of the MAC services to multiple independent users of a Bridged Local Area Network in a manner that does not require cooperation among the users, and requires a minimum of cooperation between the users and the provider of the MAC service. For example Q-in-Q can operate as follows: two VLAN tags are added to each customer Ethernet packet. The Ethernet VLAN tag includes both a 12-bit VLAN ID and a 3-bit priority tag. The inner VLAN is customer assigned, and the outer VLAN corresponds to the carrier's assignment of a tunnel in which customer traffic is carried. For example, multiple customer flows may be “clustered” into the outer VLAN. A part of the VLAN assignment can include a 3-bit p-bit priority marking per 802.1p. The VLAN assignment can indicate how the traffic is to be prioritized. Sometimes, the Ethernet packet priority is derived from the IP layer's DiffServ Code Point (DSCP) bits. The VLAN assignment can also indicate the shaping assigned to an Ethernet flow cluster.\nIn another example, PBT is used as a similar alternative to Q-in-Q except that two MAC addresses are used instead of two VLAN IDs. Each of these MAC addresses can be associated to a VLAN also, and as such has the priority markings. This is similar to the Q-in-Q except that the tunnel label field includes the MAC and that two MAC addresses are used instead of two VLANs. Another difference with PBT is that the management system normally assigns the MAC/VLAN labels along the tunnel path Ethernet switches. Once again, prioritization and clustering is done via the MAC+VLAN label.\nTransport Multi-protocol Label Switching (T-MPLS) is an ITU-T defined network layer technology that uses a subset of the existing MPLS standards and is designed specifically for application in transport networks. T-MPLS offers a simpler implementation by removing features that are not relevant to connection-oriented packet-switched applications and adding mechanisms that provide support of critical transport functionality. For prioritization in T-MPLS, an MPLS shim header is used. In this scheme, the MPLS label is used to define a traffic engineered path for with the Ethernet frames must follow. The QoS can be explicit or implicit depending on the label implementation/policy.\nDisadvantageously, all of the above described methods require prioritization to be marked in either the form of VLAN bits or MPLS labels, i.e. pre-determined or explicit. This leaves little flexibility in the grouping of individual Ethernet flows for the purpose of shaping and prioritization. For example, all Ethernet packets in one VLAN are treated with the priority of that VLAN."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for the selective growth of an epitaxial layer on source/drain areas.\n2. Description of the Related Technology\nIn semiconductor fabrication a hard mask (HM) is applied in order to protect part of the wafer from etching the other part of the wafer. The hardmask usually has to be removed afterwards. Typically an etch resistant material is employed for the hardmask.\nWhen considering a process flow for fabricating a CMOS FET transistor, a gate stack is formed, which is capped with a silicon oxide (2) and with spacers (4)(see FIG. 1(A)). Spacers (4) around the gate are typically in Si nitride. The full encapsulation is of vital importance, as it protects the gate during epitaxial growth. Further processing steps are depicted in FIG. 1. In a first step either a single layer hardmask (5) is deposited (FIG. 1(B.1)) or a hardmask stack, which cannot be etched during source/drain etching (in FIG. 1(B.2) a two-layer stack (5, 6) is shown). Next a first area is opened up comprising a source and drain (S/D) and a gate stack (FIG. 1(C)). The hardmask (5) is etched selectively towards Si oxide (used to cap the gate) and Si nitride (used for the spacers) of the first area and also selectively to the hardmask on top of a second area of the substrate (1). If the hardmask etching is not performed in a sufficiently selective way, an accurate control of the lost amount of Si oxide/Si nitride is required. In the following step (FIG. 1(D)) a source/drain (S/D) silicon etch is performed selectively towards Si oxide and Si nitride. The etch rate depends on the amount of exposed Si. A part of the hardmask (5) is consumed by a break through (BT) etch. Next an epitaxial layer is grown. During the epitaxial growth a selective growth towards the hardmask in the second area is required. The epitaxial growth is only selective towards Si oxide and Si nitride. Finally the hardmask is removed on the second area (FIG. 1(E)). Note that the hardmask cannot be removed before the epitaxial layer has grown as it should protect (at least the second area on) the wafer. In the state of the art approach the hardmask on the second area is removed afterwards by an etching step, which causes field oxide loss. A possible solution exists in providing an additional protective mask on the first area, which clearly would increase the complexity.\nSilicon nitride is the typical material for spacers (4). Si nitride cannot be applied for the hardmask (5), as during hardmask removal it would remove the spacers (4) as well. An alternative would be the use of Si oxide. However, Si oxide is not an appropriate material either, as Si oxide may be used on top of the gate. When using a Si oxide hardmask, one would have to remove the hardmask very carefully in order not to damage the oxide capping layer on top of the gate (in order to protect the gate during the epitaxial growth). Using a Si oxide hardmask, the extra oxide on top of the second area must be removed, which increases field loss on the first area, as the latter ‘sees’ the double removal of the oxide layer. An option would be to use lithography and etch only on the second area, which increases complexity.\nIn the state of the art solutions it is neither possible to use poly-silicon as the hardmask material. The poly-Si hardmask would have to be removed before the epitaxial growth, as it does not allow such epitaxial growth in a selective way. The epitaxial layer will cover the whole wafer and not only fill the S/D areas of the first area. However, it would be desirable and advantageous to have a poly-Si hardmask, as it would allow etching selectively towards nitride and oxide (e.g. towards an oxide gate capping layer or towards the field oxide), which would put less strain.\nIt is thus desirable to provide a method for selective epitaxial growth of S/D areas that overcomes the drawbacks of the prior art solutions."}
{"text": "Electrical cables are commonly supplied to an electrical junction box for connecting the cable to electrical fixtures or other electrical cables. The cables generally pass through an opening in the wall of the electrical box and are secured in place by a suitable clamp. Conventional cable clamps generally clamp the cable between the clamping plate and a side wall or bottom wall of the electrical box. This type of clamping arrangement generally relies on the smooth surface of the wall of the electrical box that does not adequately grip the cable.\nCable clamps are produced which are able to accommodate metal armor cable and plastic sheath cables. These devices include a stop member for the metal sheath to prevent the metal sheath from extending into the electrical box. Various clamping devices are known for clamping the electrical cable. Many of these cable clamps are mounted to the cable and then inserted into the opening in the electrical box so that at least a portion of the clamp projects outwardly from the outer surface of the electrical box.\nOne example of an electrical cable clamp is disclosed in U.S. Pat. No. 4,316,999 to Nattel which discloses an electrical wiring box and cable clamp where the cable clamp is attached to the wall of the electrical box. The cable clamp is a flexible member that overlies the opening in the electrical box bottom edge of the leg portions including upwardly turned flanges. The bottom wall of the electrical box includes ridges or score lines for engaging the nonmetallic cable.\nU.S. Pat. No. 4,922,057 to Ross discloses an electrical box having a flexible flat retaining member attached to the inner face of the side wall of the electrical box. As shown in FIG. 2, the flat retaining member flexes to grip the surface of the cable passing through the opening in the electrical box.\nU.S. Pat. No. 4,724,282 to Troder discloses a nonmetallic electrical box having a flexible metal cable retaining member attached to the inner face of the electrical box. The retaining member flexes when the cable is passed through the opening to grip the cable.\nU.S. Pat. No. 4,277,641 to Bauer et al. discloses a nonmetallic electrical box having a cable clamp member that is inserted into the cavity of the electrical box. The cable clamp member includes an outer frame and a flexible flap. The flexible flaps bend upon insertion of the electrical cable through the opening in the frame to grip the cable.\nU.S. Patent Publication No. 2012/0024596 to DiLillo et al. discloses a plastic cable clamp adapted for mounting in one of the openings in the electrical box. The cable clamp includes one or more flexible members capable of clamping the nonmetallic cable. In one embodiment, the cable clamp includes flexible legs that can bend inwardly upon insertion of the cable to grip the outer surface of the cable.\nWhile the prior cable clamps and connectors are generally suitable for the intended use, there is a continuing need in the industry for improved cable clamps and connectors."}
{"text": "With the advent of the Internet has come the opportunity for the print media industry to publish “on-line” rather than producing and distributing traditional paper based publications. On-line publishing has a number of advantages. From the consumer's point of view, information is available on demand, information can be navigated via hypertext links, information can be searched, and information can be automatically personalized. From the publisher's point of view, the costs of printing and physical distribution are eliminated and the publication becomes more attractive to advertisers as it can be targeted to specific demographics and linked to product sites.\nOn-line publication also has disadvantages. Computer screens are inferior to paper. At the same quality as a magazine page, a SVGA computer screen displays only about a fifth as much information. Both CRT's (Cathode Ray Tubes) and LCD's (Liquid Crystal Displays) have brightness and contrast problems, particularly when ambient light is strong. Ink on paper, being reflective rather than emissive is both bright and sharp in ambient light. Accordingly, people have a natural preference to read newspapers, magazines, catalogues, brochures and other publications in a bound paper based format.\nKnown digital printers, and in particular digital inkjet printers, do not make any provision for binding together the pages of printed media. Such an arrangement would conveniently provide for bound document production in the small office/home office (SOHO) environment. This would also permit some of the existing drawbacks of on-line publication to be addressed. However, in the SOHO environment space is often at a premium. Therefore, any provision for automatically binding separate pages into a bound document must not add a significant amount of extra space over that occupied by the printer."}
{"text": "The present invention relates to implantable systems that include medical devices (e.g., stents, vascular grafts, stent grafts) that function as a carrier for eukaryotic cells (e.g., genetically engineered endothelial cells). Such cells are capable of producing and releasing a therapeutic agent (e.g., tissue-type Plasminogen Activator) for on-demand localized treatment of conditions such as coronary artery disease. The cells or the device carrying them release the therapeutic agent upon the application of a stimulus (e.g., electrical stimulus).\nCoronary Artery Disease (CAD) affects 1.5 million people in the USA annually. About 10% of these patients die within the first year and the rest suffer from myocardial infarction and develop related symptoms, such as arrhythmias, CHF, and mechanical complications (e.g., aneursym, thrombus formation, pericarditis). During CAD, formation of plaques under the endothelial tissue narrows the lumen of the coronary artery and increases its resistance to blood flow, thereby reducing the O2 supply. Injury to the myocardium (i.e., the middle and thickest layer of the heart wall, composed of cardiac muscle) fed by the coronary artery begins to become irreversible within 0.5-1.5 hours and is complete after 6-12 hours, resulting in a condition called myocardial infarction. If the ischemia due to stenosis of the coronary artery lumen could be reduced by increasing the blood circulation to the myocardium, the major cause of most of the heart disease would be eliminated.\nCurrent and proposed treatments for coronary artery disease typically focus on pharmacological approaches and surgical intervention. For example, angioplasty, with and without stents, is a well known technique for reducing stenosis. Systemically administered drugs (e.g., anticoagulants) are also commonly used, however, such drugs become diluted, which can reduce their potency by the time they reach the remote site. Furthermore, systemic administration can be deleterious because it can lead to complications as a result of the high dosages required upon administration to allow for dilution that occurs during transport to the remote site. Therefore, localized delivery of therapeutic agents is preferred. Local delivery is advantageous in that the effective local concentration of delivered drug is much higher than can normally be achieved by systemic administration.\nStents have been used as delivery vehicles for therapeutic agents (i.e., drugs). Intravascular stents are generally permanently implanted in coronary or peripheral vessels. Stent designs include those of U.S. Pat. Nos. 4,733,655 (Palmaz), 4,800,882 (Gianturco), or 4,886,062 (Wiktor). Such designs include both metal and polymeric stents, as well as self-expanding and balloon-expandable stents. Stents are also used to deliver a drug (e.g., antiplatelet agents, anticoagulant agents, antimicrobial agents, antimetabolic agents) at the site of contact with the vasculature, as disclosed in U.S. Pat. No. 5,102,417 (Palmaz) and in International Patent Application Nos. WO 91/12779 (Medtronic, Inc.) and WO 90/13332 (Cedars-Sanai Medical Center), for example. Anticoagulant substances such as heparin and thrombolytic agents have also been incorporated into a stent, as disclosed, for example, in U.S. Pat. Nos. 5,419,760 (Narciso, Jr.) and U.S. Pat. No. 5,429,634 (Narciso, Jr.). Stents have also been used to deliver viruses to the wall of a lumen for gene delivery, as disclosed in U.S. patent application Ser. No. 08/746,404, filed Nov. 8, 1996 (Donovan et al.).\nStents seeded with autologous endothelial cells (Dichek, et al., Circulation, 80, 1347-1353 (1989)) are disclosed as a method for delivering recombinant protein over time to the vascular wall. The concentration of secreted protein produced by the endothelial cells that is required to treat the surrounding vascular tissue can be significantly higher than could be tolerated if delivered systemically. However, in order to be effective, it is not only necessary to release a high dose of the drug, but it is also necessary to achieve controlled release, e.g., immediately after an occlusion. Thus, it would be desirable to control the release of such cellular components into the surrounding tissue when needed (i.e., on demand). The present invention provides such a system.\nMany of the following lists of patents and nonpatent documents disclose information related to the local delivery of therapeutic agents using medical devices, such as stents. Others are directed toward stent designs and other medical devices as well as genetically engineered cells, for example.\nAll patent and nonpatent documents listed in Table 1 are hereby incorporated by reference herein in their respective entireties. As those of ordinary skill in the art will appreciate upon reading the Summary of the Invention, Detailed Description of Preferred Embodiments, and Claims set forth below, many of the systems, devices, and methods disclosed in these documents may be modified advantageously by using the teachings of the present invention.\nThe present invention provides implantable systems that include delivery devices having a carrier and eukaryotic cells associated therewith, optionally within a containment vehicle. Such cells are capable of producing at least one therapeutic agent (i.e., drug), which is released from the delivery device upon the application of a stimulus (e.g., electrical stimulus) for on-demand (i.e., when needed) localized treatment of conditions such as coronary artery disease or cerebral vascular occlusion, for example. The cells are referred to herein as xe2x80x9cdrug-elutingxe2x80x9d cells.\nRelease of the therapeutic agent from the delivery device is stimulated by a variety of methods, including electrical stimulation, which can be used to create an electrical field or mechanical stimulus, for example. This can be accomplished by direct action on the cells, such as by stimulating the cellular membrane to release the cellular products contained therein. Alternatively, this can be accomplished by activating the cellular products such that upon release they will function as therapeutic agents. This can also be accomplished by stimulating a microscopic containment vehicle that contains the cells, such as by stimulating a synthetic membrane of the containment vehicle, to release the cells and/or their cellular products.\nThus, an object of the present invention is to provide a system and method for the treatment (including prevention) of coronary artery disease, for example, by producing and delivering locally a therapeutic agent, such as an anticoagulant. Significantly, the local dosage can be controlled and provided on demand without worrying about the systemic effects. Furthermore, the local dosage can be administered prior to significant physiological damage occurs to the patient.\nIn a preferred embodiment, the present invention provides an implantable system that includes: a delivery device comprising a carrier for carrying eukaryotic cells that produce at least one therapeutic agent; a stimulating element for stimulating the release of the therapeutic agent from the delivery device; and a sensing element for monitoring at least one physiological property of a patient in which the system is implanted and communicating with the stimulating element to stimulate the release of the therapeutic agent from the delivery device.\nIn another embodiment, an implantable system includes a delivery device comprising an intraluminal stent, which includes a lumen-wall contacting surface, a lumen-exposed surface, a first polymer composition covering at least a portion of the stent (preferably, at least a portion of the lumen-exposed and the lumen-wall contacting surfaces), and endothelial cells associated with the first polymer composition covering, wherein the endothelial cells produce at least one therapeutic agent. The implantable system further includes an electrical stimulating element for stimulating the release of the therapeutic agent from the delivery device, and a sensing element for monitoring at least one physiological property of a patient in which the system is implanted and communicating with the electrical stimulating element to stimulate the release of the therapeutic agent from the delivery device.\nThe present invention also provides a method of local delivery (as opposed to systemic delivery) of a therapeutic agent. The method involves simply implanting an implantable system described above. Once implanted into a patient, the therapeutic agent is released from the delivery device when the sensing element detects a predetermined level of a physiological property (e.g., a certain pH or level of blood flow or blood gases) and communicating with the stimulating element to trigger the release.\nA method of making an implantable system for local delivery of a therapeutic agent is also included within the scope of the present invention. The method includes: providing a delivery device comprising a carrier and eukaryotic cells associated therewith that produce at least one therapeutic agent; providing a stimulating element for stimulating the release of the therapeutic agent from the delivery device; providing a sensing element for monitoring at least one physiological property of a patient in which the system is implanted; and connecting the stimulating element and sensing element such that they communicate with each other to stimulate the release of the therapeutic agent from the delivery device when implanted in the body of a patient and the sensing element detects a predetermined level of a physiological property. Preferably, the step of providing a delivery device includes: providing a carrier; isolating eukaryotic cells from a patient; culturing the eukaryotic cells; delivering nucleic acid of a desired therapeutic agent to the eukaryotic cells to form genetically engineered eukaryotic cells; and contacting the carrier with the genetically engineered eukaryotic cells.\nIn an alternative embodiment, the present invention provides implantable systems as described above without the sensing element. In such systems, the patient in which the system is implanted activates the stimulating element when desired (i.e., on demand). Typically, this occurs when the patient detects a change in a physiological condition (e.g., angina) and communicates with the electrical stimulating element to trigger release of the therapeutic agent. The patient can communicate with the electrical stimulating element to activate it using radio frequency, infrared, or acoustic pulses."}
{"text": "1) Field of the Invention\nThis invention relates to new antibacterial spirocyclic 6-amido-carbapenems and to a process for their synthesis through new intermediates.\n2) Brief Description of Disclosures in the Art\nCarbapenem antibiotics, particularly imipenem, (see U.S. Pat. Nos. 3,950,377 and 4,194,047) are well known for treating a broad spectrum of gram negative and gram positive bacterial infections. ##STR2##\n1-.beta.-Methylcarbapenems, as described in the reference Heterocycles, 1984, Vol. 21, pp. 29-40 by D. H. Shih, F. Baker, L. Cama and B. G. Christensen, are extremely useful and effective broad spectrum antibiotics, useful against a wide variety of bacteria including gram-positive bacteria such as S. aureus, Strep. sp., B. subtilis, and gram negative bacteria such as E. coli, Shigella sp., Enterobacter sp., Klebsiella sp., Proteus, Serratia and Pseudomonas sp.\nHowever, all of the above antibacterial carbapenems utilize 6-substituents other than amido or substituted amido which are the 6-substituents of choice in penicillin, e.g. ##STR3## or the cephalosporins, e.g. ##STR4##\n6-Amidocarbapenems and penams are known in the art as exemplified in the following references: U.S. Pat. No. 4,260,627; U.S. Pat. No. 4,206,219; U.S. Pat. No. 4,217,453; U.S. Pat. No. 4,218,459; U.S. Pat. No. 4,218,463; U.S. Pat. No. 4,277,482; and U.S. Pat. No. 4,298,741 to Merck & Co., Inc. which describe 1H-6-amidocarbapenems and 1-methyl 6-aminocarbapenems; BE 887,618 and U.S. Pat. No. 4,347,355 to Abbott which describe 1,1-diloweralkyl-6 -amidocarbapenems; EPO Publication No. 040,494 and U.S. Pat. No. 4,348,264 to Pfizer which describes 1-hydroxy, acetoxy or 1,1-oxo-carbapenems with 6-position conventional penicillin sidechains; EPO Publication Nos. 634,443 and 073,100 and U.S. Pat. No. 4,407,815 to Beecham which describe 1-H-6-amidocarbapenems and penams; Japanese Kokai 58 174 382 to Sanraku Ocean Co. Ltd. which discloses 6-phthalimido-2-SR carbapenems; and EPO Publication No. 045,198 to Takeda Chem. Ind. Ltd. which discloses 1-alkyl-1-alkoxycarbonyl, cyano or COR-substituted-6-aminocarbapenems.\nLiterature articles relating to 1-H-6-amidocarbapenems discussing problems in ring closures and ester deblocking reactions include: Tetrahedron Letters, 1982, 23 (15), 1545-1548; Tetrahedron Letters, 1982, 23 (50), 5339-5342; L. C. Blaszczak, Eli Lilly Co. Report \"Joint Great Lakes and Central Regional Meeting\", Western Michigan University, May 23-24, 1984; J. Chem. Soc., Perkins Trans. I, 1982, 2123-2129; Helv. Chim. Acta. 1982, 65, 1378-1384; J.A.C.S., 1982, 104, 4262-4264; Chem. Pharm. Bull. 31 2578 (1983); N. Narisada et. al. 176th ACS National Meeting, Miami, Fla. 1978.; and Nouv. J. Chim. 7 691 (1983).\nThe assignee herein has been responsible for more recent work with 6-amidocarbapenems. U.S. Ser. No. 213,579, filed Jun. 30, 1988, discloses novel 6-amido-1-methyl-2-(substituted thio)-carabapenems. Also, U.S. Ser. No. 213,398, filed Jun. 30, 1988, discloses novel 6-amido-1-methyl-carbapenems. Further, assignees docket No. 17660 filed the same date herewith discloses novel 6-amido-1-methylcarbapenems. Each class of compounds exhibits antibacterial activity.\nNew antibacterial compounds are constantly being searched for to enhance the potency and decrease the side effects of current existing carbapenem antibiotics. Thus far, spirocyclic 6-amido carbapenems have not been disclosed in the art."}
{"text": "Floor care appliances such as upright vacuum cleaners have long been known in the art to be useful in cleaning dirt and debris from floors and the like. Generally such vacuum cleaners are not suited for use by individuals with disabilities that restrict them to movement in a wheelchair. The present invention relates to a vacuum cleaner particularly adapted for use by such individuals."}
{"text": "When making a housing for an electronic device for use in outdoor environments, an important consideration is that the housing should be able to cope with large differences in temperature. That is, the housing should be adapted to both cold and hot environments. Much focus is hence needed on insulation and power consumption, and from a economical and environmental perspective it is desirable to keep power consumption at a low level.\nProtective housings for electronic devices for outdoor environments usually have a single sealed shell for the purpose of protecting the electronic devices inside the shell. Heat in the shell generated by the electronic device can easily be conducted through the shell if the external air temperature is low. However, when the external temperature is high or the shell is heated by solar radiation, the internal temperature will rise due to the generated heat from the electronic device and the external heat. A solution to such a problem may be to arrange a sunshade or adding a shell outside the sealed shell.\nConventional double shell housings for electronic devices use air between the shells to prevent the external heat to reach into the inner shell. A ventilation opening is typically provided in the outer shell, and provide a natural convection of the air if the inside of the inner shell is heated by the electronic device to a higher temperature than the external temperature. Another way is to manufacture shells in a heat conducting material so as to facilitate transmission of heat, which is generated by the electronic device, to the outside in order to keep the internal temperature down.\nLikewise, in conventional double shell housings, the air between the shells may be useful in cold environments, as the air between the shells is applied as an insulation to maintain a convenient tempered environment inside the inner shell around the electronic device.\nIn such double shell housings using natural convection and an opening in the outer shell, however, there is a problem with a non-effective heat transportation both from the inside of the inner shell and also between the inner and outer shell.\nEP 0746 192 A1 addresses this problem by using a plurality of fans provided between an inner and an outer shell, where the inner shell is hermetically sealed and the outer shell has bimetal valves for support of circulation with the environment outside the outer shell. However, EP 0746 192 A1 has a drawback, in that the fans and any electronics connected to the fans will be exposed to the outside environment. If the outside environment is cold and damp, this will lead to damages to the fans and any connected electronics.\nIn U.S. Pat. No. 4,984,089 an outdoor surveillance apparatus is disclosed which support movement of a camera assembly. The apparatus is adapted to low and high temperature by use of a heater and multiple fans which are regulated by thermostats.\nU.S. Pat. No. 4,419,716 disclose a sealed vapor proof housing with a system for isolating an electrical device from ambient atmosphere. Heat is transferred by use of a fan and that the housing made of metal consists of a multiple number of heat exchange fins."}
{"text": "In recent years, machine-learning approaches for data analysis have been widely explored for recognizing patterns which, in turn, allow extraction of significant information contained within a large data set that may also include data consisting of nothing more than irrelevant detail. Learning machines comprise algorithms that may be trained to generalize using data with known outcomes. Trained learning machine algorithms may then be applied to predict the outcome in cases of unknown outcome, i.e., to classify the data according to learned patterns. Machine-learning approaches, which include neural networks, hidden Markov models, belief networks and kernel-based classifiers such as support vector machines, are ideally suited for domains characterized by the existence of large amounts of data, noisy patterns and the absence of general theories. Support vector machines are disclosed in U.S. Pat. Nos. 6,128,608 and 6,157,921, both of which are assigned to the assignee of the present application and are incorporated herein by reference.\nMany successful approaches to pattern classification, regression, clustering, and novelty detection problems rely on kernels for determining the similarity of a pair of patterns. These kernels are usually defined for patterns that can be represented as a vector of real numbers. For example, the linear kernels, radial basis function kernels, and polynomial kernels all measure the similarity of a pair of real vectors. Such kernels are appropriate when the patterns are best represented in this way, as a sequence of real numbers. The choice of a kernel corresponds to the choice of representation of the data in a feature space. In many applications, the patterns have a greater degree of structure. This structure can be exploited to improve the performance of the learning system. Examples of the types of structured data that commonly occur in machine learning applications are strings, such as DNA sequences, and documents; trees, such as parse trees used in natural language processing; graphs, such as web sites or chemical molecules; signals, such as ECG signals and microarray expression profiles; spectra; images; spatio-temporal data; and relational data, among others.\nFor structural objects, kernels methods are often applied by first finding a mapping from the structured objects to a vector of real numbers. In one embodiment of the kernel selection method, the invention described herein provides an alternative approach which may be applied to the selection of kernels which may be used for structured objects.\nMany problems in bioinformatics, chemistry and other industrial processes involve the measurement of some features of samples by means of an apparatus whose operation is subject to fluctuations, influenced for instance by the details of the preparation of the measurement, or by environmental conditions such as temperature. For example, analytical instruments that rely on scintillation crystals for detecting radiation are known to be temperature sensitive, with additional noise or signal drift occurring if the crystals are not maintained within a specified temperature range. Data recorded using such measurement devices are subject to problems in subsequent processing, which can include machine learning methods. Therefore, it is desirable to provide automated ways of dealing with such data.\nIn certain classification tasks, there is a priori knowledge about the invariances related to the task. For instance, in image classification, it is known that the label of a given image should not change after a small translation or rotation. In a second embodiment of the method for selecting kernels for kernel machines, to improve performance, prior knowledge such as known transformation invariances or noise distribution is incorporated in the kernels. A technique is disclosed in which noise is represented as vectors in a feature space. The noise vectors are taken into account in constructing invariant kernel classifiers to obtain a decision function that is invariant with respect to the noise."}
{"text": "This invention relates to fluid flow rate measurements. More specifically, this invention relates to determining fluid flow rates using measured motor speed and torque parameters.\nMeasurement of fluid flow through pipes and pumps is well known in the art. One method of determining fluid flow rates is to install gears, vanes, paddle wheel, turbines, etc., in the flow channel and determine fluid flow rate by the speed at which these devices turn. A second method is by measuring the differential pressure across a dedicated flow obstruction, such as a venturi, orifice plate, annubar, pitot tube, etc., and applying the well-known Bernoulli\"\"s principle to obtain a velocity and, consequentially, a fluid flow rate. As an example of this principle, U.S. Pat. No. 5,129,264, entitled xe2x80x9cCentrifugal Pump with Flow Measurement,xe2x80x9d issued Jul. 14, 1992, to Lorenc, and assigned to the same assignee herein discloses using differential pressure and pump speed to measure fluid flow rates. Still other methods of fluid flow rate measurement employ electrical/magnetic or sonic measurement means. For example, Mag Meters determine fluid flow rates by measuring the change in a magnetic field caused by the velocity of the fluid flowing therethrough. Sonic devices use acoustical pulses, i.e., Sonar, and Doppler principles to measure fluid flow rates. Other non-intrusive methods measure the torque a variable speed electrical motor delivers to a pump by installing a torque shaft between the motor and pump. The motor or pump is then calibrated and a motor kilowatt input/Motor Brake Horsepower Output (BHP) calibration table is developed. Accordingly, knowledge of the kilowatt input value can be used to determine the output horsepower. However, this calibration is needed at several speeds and requires several different sized torque shafts.\nThus, current methods for determining fluid flow rates necessitate intrusion into, or require access to, the enclosures transporting the fluid. In some systems, such as caustic systems having lined pumps, intrusion is prohibited.\nHence, there is a need for a simple, accurate and reliable method for determining fluid flow rates without intruding into the fluid flow or having access to the enclosures transporting the fluid.\nThe present invention determines a fluid flow rate through a pump by first determining two flow rate values from a plurality of characterizing flow rate values corresponding to two known speed values selected from a plurality of known characterizing speed values at a known motor torque. The first one of the two known speed values is selected greater than a measured pump speed and a second one of the two known speed values is selected less than the measured pump speed. The fluid flow rate is then determined as being proportional to the two determined flow values at the known speed values and the pump speed."}
{"text": "1. Field of the Invention\nThe present invention relates to the game apparatus, and more particularly to the game apparatus with attachable game pieces played on the outer surfaces of the three-dimensional game structure.\n2. Description of Related Art\nVarious board-type racing games have been developed in the past. U.S. Pat. No. 6,883,803 B1 issued to Barry on Apr. 26, 2005 discloses a two-dimensional tourist game board. U.S. Pat. No. 4,182,516 issued to Gill on Jan. 8, 1980 discloses a two-dimensional sailboat racing game. U.S. Pat. No. 3,871,656 issued to Selness on Mar. 18, 1975 discloses a two-dimensional sailing game apparatus. These games are two-dimensional planer games and did not suggest three-dimensional movement for game pieces on the three-dimensional game structure. The present invention is to produce a game in which the player can be challenged to race and capture the opponent's game piece on the transformable three-dimensional game structure.\nThe three-dimensional game apparatus which incorporate the concept of tic-tac-toe or checker type of games have been developed in the prior art wherein a series of two-dimensional playing surfaces are vertically arranged one above the other. So that game piece can be moved or arranged on a two-dimensional playing surface as well as between the playing surfaces. One such game apparatus is disclosed by Mahoney U.S. Pat. No. 3,464,701 on Sep. 2, 1969, wherein a series of horizontally disposed playing boards are vertically supported in parallel spaced relationship by a box-shaped frame.\nU.S. Pat. No. 3,623,729 issued to Wetherell on Nov. 30, 1971 discloses a three-dimensional board game apparatus. U.S. Pat. No. 3,656,755 issued to Thompson on Apr. 18, 1972 discloses a three-dimensional checker game apparatus. U.S. Pat. No. 4,129,303 issued to Flagg on Dec. 12, 1978 discloses a cubic game board. Another three-dimensional game apparatus of the tic-tac-toe type is constructed with rods joined together to form a grid pattern in the shape of a cube. A three-dimensional tic-tac-toe game is played by arranging game pieces on selected horizontal runs of the grid pattern. A game apparatus of this type is disclosed by Green in U.S. Pat. No. 3,606,333 on Sep. 20, 1971.\nEven though the prior art's innovations may be suitable for the specific individual purposes to which they address, they would not be comparable for the purposes of the present invention as heretofore described."}
{"text": "1. Field of the Invention\nThe present invention relates to a paint bucket comprised of an open-top-end container having a lid with attachment brackets disposed on a rear face of the lid to retain paint applicating implements thereagainst and wherein the paint bucket has a slot in a top edge thereof for the passage of a paint roller connecting rod whereby a wet roller may be stored inside the container to prevent drying thereof.\n2. Description of Prior Art\nIt is known in the prior art to provide paint buckets having a ribbed wall surface formed with the bucket whereby to remove excess paint from a roll applicator brush rolled thereagainst. This is disclosed in U.S. Pat. No. 2,988,767. However, with such buckets there is a need to remove paint from the bucket in order to use the ribbed side wall as the paint in the bucket, when full, conceals this ribbed side wall. It is therefore necessary to have another bucket to store paint while the bucket with the flat ribbed wall is utilized. U.S. Pat. No. 2,983,938 also discloses a paint bucket wherein a roller-type paint applicator is inserted within the container, which is filled with paint, through a small cylindrical orifice and then withdrawn therefrom and during this withdrawing process a plurality of circumferential serrations will be in doctoring contact with the peripheral surface of the roller brush to remove excess paint. The purpose of this device is to eliminate the need to use paint pans which need to be constantly refilled due to their small capacity size. It also eliminates the need to clean the paint pan.\nThere is, however, another need for storing wet paint applicating implements, such as rollers or paint brushes, to prevent them from becoming dry when there is a pause in the applicating of paint for a substantial period of time. When wet paint implements are left uncleaned and exposed to air they dry out. This need to prevent wet paint implements from drying also eliminates the need to have to clean the paint implements each time there is a lengthy pause. Often paint implements have to be left wet for a substantial length of time.\nThere is a further need to provide a paint bucket which will also serve as a paint tray but without having to remove paint from within the bucket for access to a ribbed surface of the paint bucket whereby excess paint from a roller-type paint applicator can be removed."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical output control circuit of a semiconductor laser for constantly setting the optical output of the semiconductor laser irrespective of ambient temperature.\n2. Description of the Related Art\nFIG. 4 shows an optical output control circuit of a conventional semiconductor laser. A reference voltage Vref is applied to the base of a transistor 42 for flowing a modulating electric current to the semiconductor laser 41, and its emitter is connected to the emitter of a transistor 43 having a base to which a pulse signal for modulation is inputted. These two emitters are connected to the collector of a transistor 44, and an operational amplifier 45 is arranged between the base and the emitter of the transistor 44. The voltage of an emitter resistor 46 of the transistor 44 is inputted to the operational amplifier 45. Further, negative feedback is applied to an inversion input terminal (−) of the operational amplifier 45 by a thermistor 47. A non-inversion input terminal (+) of the operational amplifier 45 is connected to the ground by a thermistor 48.\nIn the above construction, the transistors 42 and 43 are alternately turned on/off by a continuous pulse signal so that the modulating electric current is flowed to the semiconductor laser 41. Further, since the thermistors 47, 48 are connected to the operational amplifier 45, the semiconductor laser 41 is operated such that the modulating electric current flowed to the semiconductor laser 41 is increased correspondingly with a rise in temperature. Thus, the optical output is constantly set even when the semiconductor laser 41 is reduced in efficiency.\nFurther, one portion of the optical output of the semiconductor laser 41 is detected by a light receiving element 50. The detected voltage is inputted to an optical output control circuit 51. The optical output control circuit 51 operates a transistor 52 for flowing a bias electric current to the semiconductor laser 41 such that the optical output of the semiconductor laser 41 is constant irrespective of the ambient temperature. Namely, if the detected voltage is reduced by a reduction in efficiency of the semiconductor laser 41 caused by the rise in temperature, the optical output control circuit 51 controls the operation of the transistor 52 so as to compensate for this reduction in the detected voltage so that the bias electric current is increased (see Patent Literature 1).\n[Patent Literature 1]\nJapanese Unexamined Published Patent Application No. H5-259563 (FIG. 1)\nThe level of the optical output outputted from the semiconductor laser 41 is determined by a sum of the bias electric current and the modulating electric current. Therefore, if the pulse signal for modulation is continuously inputted in the above construction, the modulating electric current and the bias electric current are changed dependently on temperature. Accordingly, the optical output level is controlled by the optical output control circuit 51 such that the optical output level is constant at all times. However, when the pulse signal for modulation is inputted in a so-called burst shape, no modulating electric current is flowed to the semiconductor laser 41 in a period in which no pulse signal for modulation is inputted. Accordingly, no electric current is flowed to the light receiving element 50 so that the optical output control circuit 51 increases the flow electric current to the semiconductor laser 41 to a maximum. As this result, there is a problem of destroying the semiconductor laser 41. Accordingly, no semiconductor laser can be used in the burst mode in the conventional construction."}
{"text": "Not applicable.\nThe present invention relates to a marine or other recreational vehicle receptacle for receiving and storing a marine shower. The marine shower may or may not accommodate a hot-cold water mixing valve as required by the consumer.\nWhile the invention has multiple uses readily apparent to those skilled in the art it will be described primarily for use in connection with the marine industry.\nMarine showers are most often mounted in the transom of a boat or in the cockpit coaming of a boat and in modern day boat designs these transoms and coamings comprise a multitude of surfaces with orientations from vertical to horizontal and therebetween. These orientations require different shower head receptacle designs and when a mixing valve is required a still further design is needed. Thus multiple items must be kept in inventory in order to accommodate all requirements.\nA shower receptacle designed for a vertical surface will not be useful for a horizontal surface and a receptacle adapted for a straight shower handle will not accommodate a Euro-style handle with a projecting on-off lever. When a hot-cold water mixing valve is desired a further receptacle is needed to accommodate the valve. Some prior art designs feature the shower wand and hot and cold water valves in one box but this requires a large opening to be cut in the boat transom with its attendant disadvantages.\nThe disadvantages and limitations of the prior art devices are obviated by the present invention.\nAn object of the present invention is to provide a shower head receptacle adapted to be mounted on surfaces having different angles of orientation.\nA further object of the present invention is to provide a shower head housing or receptacle for receiving either a straight or Euro-style shower wand or handle.\nA still further object of the present invention is to provide a receptacle with a lid which is relatively easy to manufacture.\nAnother object of the present invention is to provide a receptacle which can accommodate various shower heads or wands as well as a hot and cold water mixing valve as the case may be.\nAnother object of the invention is to provide a receptacle for a shower head and a receptacle for a mixing valve in side by side relation with a single cover for both receptacles.\nOther objects, advantages and features of the present invention will become apparent to those skilled in the art from the following detailed description which with referencess to the accompanying drawings discloses a preferred embodiment of the invention."}
{"text": "The present invention relates to steam turbines and more particularly to a steam turbine cover for minimizing or eliminating solid particle erosion of bucket tenons and to improve sealing efficiency.\nIn conventional steam turbines, bucket covers are typically attached to the buckets by peening the tenon projecting from the end of the bucket and through an aperture in the bucket cover. This peening operation results in a projecting bulb or knob on the outside diameter of the cover. This raised knob or projection can be eroded by solid particles in the steam path. As a result, the cover may become loose, crack or separate from the buckets. Also, the raised knob or projection, resulting from peening the tenon, substantially prevents the application of one or more labyrinth seal teeth along the outside face of the cover, particularly on units with large differential expansion between the stationary and rotating components.\nIn another cover-to-bucket design, the knob or projection of the peened tenon is disposed below the outer surface or outer diameter of the cover. While that configuration enabled the application of sealing teeth configurations to the cover, it is limited to straight tooth or flat surface geometry. Under certain conditions, solid particles may become trapped within the confines of the recess and between the walls defining the recess and the tenon. These trapped particles tend to erode the tenon knob or projection very quickly and have been shown, in certain circumstances, to essentially cut through the tenon.\nAnother bucket/cover design includes the integral formation of the bucket and cover. While this design incorporates the necessary sealing options, i.e., application of one or more labyrinth seal teeth, and also minimizes or eliminates the concern for solid particle erosion, the integral bucket/cover combination is costly to manufacture and complex. Accordingly, there has arisen a need for a bucket/tenon/cover design that both eliminates or minimizes solid particle erosion, as well as affords sealing efficiencies without complexity or excessive costs.\nIn accordance with a preferred embodiment of the present invention, there is provided in a steam turbine, a plurality of buckets mounted on a rotating component, e.g., a rotor, and a plurality of covers mounted on the tips of the buckets, the buckets having tenons peened to secure the buckets to the covers. At least one recess is preferably formed in the outer cover surface to form a profiled surface and at least one tooth projects either radially outwardly from the profiled surface or radially inwardly from the registering stationary component. In a preferred embodiment hereof, the profiled surface includes at least one recess defined at least in part by the tenon. In another preferred embodiment, the profiled surface includes a tooth projecting from the outer surface of the cover toward the stationary component and including at least part of the tenon. It will be appreciated that the recess or the tooth each formed, at least in part, by the tenon extend in a circumferential direction about the cover. To form this preferred configuration, the outer surface of the cover is machined such that the tenon forms part of the recess or the labyrinth tooth, as applicable.\nIn a preferred embodiment according to the present invention, there is provided in a steam turbine having a plurality of buckets rotatable about an axis and a stationary component surrounding the buckets, a seal between the buckets and the stationary component, comprising a cover mounted on a radial outer end of at least one bucket and having an opening, one bucket having a tenon projecting from the outer end of one bucket and extending into the cover opening, an outer surface of the cover and an outer end surface of the tenon forming a profiled surface in opposition to the stationary component with contiguous surfaces of the tenon end and the cover lying flush with one another, the profiled surface including at least one of a recess formed in the profiled surface and a tooth projecting radially outwardly of the profiled surface.\nIn a further preferred embodiment according to the present invention, there is provided in a steam turbine having a plurality of buckets rotatable about an axis and a stationary component surrounding said buckets, a seal between the buckets and the stationary component, comprising a plurality of covers mounted on radial outer ends of the buckets arranged in an annular array thereof with one or more buckets being secured to each cover, each cover having at least one opening and each bucket having a tenon projecting from the outer end thereof into the opening, an outer surface of each cover and an outer end face of each tenon forming a profiled surface in opposition to the stationary component with contiguous surfaces of the tenon end faces and the outer cover surface lying flush with one another.\nIn a further preferred embodiment according to the present invention, there is provided in a steam turbine having a plurality of buckets rotatable about an axis and a stationary component surrounding the buckets, a seal between the buckets and the stationary component, comprising a plurality of covers mounted on radial outer ends of the buckets arranged in an annular array thereof with one or more buckets being secured to each cover, each cover having at least one opening and each bucket having a tenon projecting from the outer end thereof into the opening, an outer surface of each cover and an outer end face of each tenon forming a profiled surface in opposition to the stationary component with contiguous surfaces of the tenon end faces and the outer cover surface lying flush with one another, the profiled surface including at least one recess, the profiled surface including a tooth projecting radially outwardly thereof."}
{"text": "FIG. 8 shows a conventional nozzle through which air can be pumped into a cushion (not shown) from a pump (not shown). The cushion includes a tube and a valve installed in the tube. A thread is formed on an internal face of the tube.\nThe nozzle comprises a housing 100, a tube 200 and a tubular seal 300. The housing 100 comprises an annular portion and a cylindrical portion 120 extending from the annular portion. A thread 122 is formed on the annular portion of the housing 100.\nThe tube 200 includes a first section inserted through the annular portion of the housing 100, a second section put in the housing 100 and an annular flange 210 formed between the first and second sections of the tube 200 for abutment against the annular portion of the housing 100. The first section of the tube 200 is inserted in a pipe 400 leading to a pump (not shown). A thread 212 is formed on the first section of the tube 200 for engagement with the thread 122 formed on the annular portion of the housing 100. A ferrule 420 is put around the pipe 400 in order to keep the first section of the tube 200 in the pipe 400. Pushers 220 are formed at the second section of the tube 200. The pushers 220 leave the second section of the tube 200 open.\nThe tubular seal, 300 is put around the second section of the tube 200. The tubular seal 300 defines a big hole 320 and a small hole 340 communicated with the big hole 320. A portion of a tubular lip 140 is screwed in the cylindrical portion 120 of the housing 100. A thread 142 is formed on an internal face of the tubular lip 140.\nTo pump air, a user rotates the housing 100 so as to rotate the tubular lip 140 in order to engage the thread 142 of the tubular lip 140 with the thread formed on the tube of the cushion. Now, the valve of the cushion is pushed by the pushers 220 so that air flows through the tube of the cushion. The user rotates the tube 200. Because of the engagement of the thread 212 with the thread 122, the tube 200 is moved towards the tube of the cushion. The tubular seal 300 will also be moved to the tube of the cushion and seal the same. The pressurized air can be pumped into the cushion from the pump. However, much air has escaped from the cushion after the valve installed in the tube of the cushion is pushed by the pushers 220 before the tube of the cushion is sealed by the tubular seal 300. In addition, it takes quite some time to rotate the tube 200 in order to use the seal 300 to seal the tube extending from the cushion.\nThe present invention is therefore intended to obviate or at least alleviate the problems encountered in the prior art."}
{"text": "Games played on tablet computers typically rely on the touch screen interface of the tablet device to interact with the user."}
{"text": "1. Field of the Invention\nThis invention relates to technology of manufacturing a semiconductor device. More especially, this invention relates to a set of masks, a method of generating mask data and a method for forming a pattern.\n2. Description of the Related Art\nIn recent years, with the trend toward higher integration of large scale integrated circuits (LSIs), there has been a demand for finer circuit patterns of LSIs. As examples, for a transistor in a microprocessing unit (MPU), an application specific integrated circuit (ASIC), a system LSI, or the like, there is a demanded for a pattern to be delineated with fine line widths of 100 nm or less.\nAs a method for forming a pattern with fine line widths, there is a known method in which a resist pattern with fine line widths is formed using a positive alternating phase shift mask (Levenson-type phase shift mask) where phase contrasts of adjacent openings are shifted from each other by 180 degrees, and any unnecessary part of the resist pattern is removed using a trim mask. Moreover, as another method for forming a pattern with fine line widths, a method has been studied in which isotropic etching is performed after the formation of a resist pattern, thereby slimming the line widths of the resist pattern.\nHere, for example, as to a pattern of a memory cell of a SRAM or the like which coexists with a logic circuit, demand for the pattern have fine space widths in addition to fine line widths in order to reduce the circuit space. In the case of using a positive alternating phase shift mask (Levenson-type phase shift mask), it is required to form fine space widths when a trim mask is used. However, the resolution of a trim mask is low, and it is therefore difficult to form fine space widths in a resist pattern. Moreover, in the case of performing isotropic etching after the formation of a resist pattern, the space widths between the lines of the resist pattern are increased when the line widths of the resist pattern are slimmed. Therefore, it is difficult to delineate a circuit pattern which has fine line widths and fine space widths at the same time."}
{"text": "In this type of installation, the ladles containing the molten metal coming from the steel plant are brought near to the casting installation by means of a transport carriage or a travelling crane and are placed above a tundish, into which the molten metal pours via a taphole. The molten metal subsequently passes from the tundish to the bottomless mold, where the continuous strip of metal, namely a slab or a bar, is formed.\nIn order to make it easier to replace the ladles and reduce the downtime as much as possible, an appliance called a \"ladle turntable\" is often used as a ladle support.\nSuch appliances are known. They comprise at least one arm extending in an overhung manner from a barrel mounted rotatably about a vertical axis on a stationary base, for example the supporting stand of the installation. Each arm can support a ladle and, as a result of rotation, move it from a charging and discharging position into a casting position above the tundish.\nThe turntable is often equipped with two diametrically opposite arms, so that, while one ladle located above the tundish pours the steel which it contains, the previously emptied ladle can be removed from the other arm in order to replace it with a new full ladle.\nIn general, the taphole closed by means of a spout is extended downwards by a casting pipe, the function of which is to protect the jet of steel between the ladle and tundish from contact with the air and therefore from oxidation.\nThe protective devices are generally either a refractory pipe attached or retained under the ladle or a skirt made of asbestos or similar material, in which an inert atmosphere is maintained.\nBecause the ladle is brought above the tundish and then emptied as a result of the rotation of the arm of the support about its vertical axis, it was thought expedient to equip the ladle supports with lifting devices which make it possible to release the casting pipe during changes of ladles or tundishes. Lifting is also necessary when, in order to unblock the spout closing the taphole, it is necessary to blow into it, by means of a lance, pressurized oxygen which burns the steel solidified in the spout orifice.\nThe known devices for lifting the ladle-supporting arms can be of the hydraulic or the electromechanical type.\nIn devices of the mechanical type, lifting is brought about by means of elements of the winch and block and tackle, screws and nuts, connecting rod or crank type. These devices are heavy and bulky, and they sweep across a large space through the rotation of the turntable, thus creating dangerous zones at the workstations. Furthermore, they generally comprise a certain number of joints and guides which are exposed to the splashes of steel, and which therefore require frequent and difficult maintenance.\nMoreover, the guides projecting above the arms require a greater lifting height on the crane and therefore a proportionate building height.\nIn devices of the hydraulic type, lifting is brought about by means of jacks. These devices are simpler and generate fewer stresses than the mechanical devices.\nIn this case, the lifting of the ladle is controlled by at least one hydraulic jack bearing on the supporting arm, the latter being mounted pivotally on the barrel about a horizontal axis.\nThe ladle is equipped with two journals engaging into suspension members which allow actuation by means of the travelling crane, and which rest on bearing parts formed at the end of the supporting arm and often having the form of a trapezoidal cup, into which fit the lower parts, of corresponding forms, of the suspension members. Thus, if the bearing parts are fastened rigidly to the end of the supporting arms, the pivoting of these during lifting determines a change of orientation of the supporting members and consequently an oscillation of the ladle about its journals. Consequently, it is often considered preferable to provide the bearing parts on two supporting stays connected to the barrel by means of an articulated system constituting a deformable parallelogram, of which they form an outer side moveable parallel to itself. For this purpose, each supporting stay is articulated about horizontal axes, respectively on the end of the supporting arm and on a connecting rod parallel to the supporting arm and articulated on the base, the vertex of the deformable parallelogram being formed by the respectively stationary and moveable axes of articulation of the supporting arm and of the retaining connecting rod respectively on the base and on the two supporting stays.\nIn such arrangements, which are described, for example, in EP-A-0,206,169 or FR-A-2,234,946, the lifting force is exerted on the supporting arm, which is therefore subjected to bending, and the force required from the jack is greater than the weight of the charge. For this reason, in order to support the articulated assembly more effectively during lifting, the latter is generally produced by means of two set-apart and synchronized jacks, each acting on one side of the supporting arm."}
{"text": "The present invention relates to a grommet constructed and arranged to make a blind connection between a primary panel and a secondary panel. More specifically, the grommet of the present invention may be laterally secured to a pin extending from the undersurface of the primary panel.\nThe present invention can best be summarized as a side mounting grommet that comprises a body having an aperture formed entirely therethrough and a head that extends laterally from its upper end and which acts to limit the depth of insertion of the body into a bore that is formed into a panel. Both the body and the head have a slot formed therethrough that communicates with the aperture formed through the body. This slot allows the grommet to be addressed laterally to a pin that extends from the underside of a primary panel such that the pin is received within the aperture.\nA detent structure formed within the aperture of the body extends into the aperture and is constructed and arranged to engage a flat or depression formed near the base of the pin such that once the pin is seated within the aperture of the grommet\"\"s body, the pin may not easily be withdrawn longitudinally from the bore. At least one retaining mechanism on the exterior surface of the body of the grommet acts to retain the grommet, and the pin that has been pre-assembled therewith, within a bore formed in a secondary panel. Typically, the connection made between the primary and secondary panels is considered xe2x80x9cblindxe2x80x9d, i.e., the pre-assembled pin and grommet are not visible as they are being inserted into the bore formed through the secondary panel.\nOne advantage of the grommet of the present invention is that it may be disassembled for the purposes of repair or replacement of the panels or related structures. To facilitate this, the cooperative bearing surfaces of the depression formed in the pin and/or the bearing surfaces of the detent structure may be constructed and arranged to impose a predetermined withdrawal force upon the grommet. Forming an angled or curvilinear surface on one or both of the detent structure and depression will allow one or both of these structures to be biased or elastically deformed such that the pin may be removed longitudinally from the aperture in the grommet body. The withdrawal force required to remove the pin longitudinally from the body may be larger or smaller than the force required to remove the body of the grommet from the bore formed in the secondary panel when the grommet and pin are received therein.\nTo further facilitate the disassembly of the grommet of the present invention, the head of the grommet may be provided with release apertures that allow a user to pry or bias the retaining mechanism(s) that retain the grommet body in the bore in the secondary panel inwardly such that the grommet body may be removed from the bore.\nThe present invention also comprises a method of creating a blind connection between a primary and a secondary panel. The process begins with the step of providing a pin on the underside of the primary panel. This pin may be formed integral with the panel in a molding operation or it may be adhered or mechanically secured thereto after the primary panel has been manufactured. A grommet having a body constructed as described herein is then assembled with the pin by addressing a longitudinal slot in the grommet body to the pin and forcing the grommet body onto the pin such that the pin is received and retained within the aperture formed through the body of the grommet. Once the grommet and pin have been pre-assembled, the grommet, with the pin received therein, is inserted into a bore formed through the secondary panel until the retaining structures on the grommet body engage the secondary panel, thereby securing the primary panel to the secondary panel. Disassembly of the primary and secondary panel is undertaken by first applying a withdrawal force to the pin. Note that this withdrawal force is typically applied by prying the primary panel away from the secondary panel but may be applied directly to the pin itself where a suitable tool can be addressed thereto. Once the pin has been withdrawn from the grommet body, the retaining structures are accessed through the release apertures formed in the head of the grommet body and biased away from the bore in the secondary panel to a degree that allows the grommet body to be removed from the bore.\nAlthough the disclosure hereof is detailed and exact to enable those skilled in the art to practice the invention, the physical embodiments herein disclosed merely exemplify the invention, which may be embodied in other specific structure. While the preferred embodiment has been described, the details may be changed without departing from the invention, which is defined by the claims."}
{"text": "Various collaborative communications platforms exist, such as unified telephony systems, that may allow users of the system to communicate with one another for various purposes. For example, such systems may support multi-party communications, such as conference calls, web-conferences, virtual meetings, and the like. As is sometimes the case, one or more parties on a call may do most of the talking, with other parties participating in a more passive or listening role. In some such circumstances, the speaking parties may not be aware of changes in circumstances of other participants on a call, which can lead to confusion, for example, if a response, or other input is solicited from generally passive or listening party on the call. For example, such a passive or listening party may have become distracted from participation in the call, but due to their passive role on the call, the other participants may not be aware of the change in circumstance."}
{"text": "When making a CMOS device that includes metal gate electrodes, it may be necessary to make the NMOS and PMOS gate electrodes from different materials. A replacement gate process may be used to form gate electrodes from different metals. In that process, a polysilicon layer, bracketed by a pair of spacers, is removed to create a trench between the spacers on both NMOS and PMOS gate areas simultaneously. The trench is filled with a first metal. Then, the first metal layer is removed from either the NMOS or PMOS gate area and replaced with a second metal that differs from the first metal.\nThe replacement gate process requires the polysilicon and metal layers to be formed on a dielectric layer. The dielectric layer serves as an etch stop layer and prevents significant numbers of ions from reaching the channel when ions are implanted into the polysilicon layers. The dielectric layer may, for example, comprise silicon dioxide or a high-k dielectric layer.\nU.S. Pat. No. 7,160,767 to Brask, et al. (“Brask”), provides a method for making a semiconductor device that includes metal gate electrodes, wherein such method includes replacing polysilicon layers with metal layers without damaging the channel region. However, Brask also has disadvantages. For example, the Brask method includes a large number of steps and is otherwise complex, such as its required removal of two different dummy gates, its deposition of two different high-K dielectric layers, and its performance of two different CMP processes. The Brask method is also disadvantageous in that gate height is lost due to the requirement for two different CMP processes. The complex and lengthy method taught by Brask is also very expensive. Brask also requires thick layers of precious metals, which are expensive to obtain. Moreover, the Brask method results in a metal gate plagued by internal seams, gaps or voids which, when not sufficiently filled, results in excessively high gate resistance."}
{"text": "This invention is concerned with a method and system for tenderizing meat. A standard technique for improving the flavor of meat is to soak meat before cooking in a marinade. When meat is placed in a liquid marinade, the flavor is enhanced due to the absorption of some of the marinate liquid into the meat. Others have suggested enhancing the effect of marination by including enzymes.\nTo speed up and to augment the process of the absorption of a marinade by meat, others have suggested placing the marinade in a closed container in which the meat is placed and subjecting the container to pressure--the concept being to utilize fluid pressure to force more of the marinade into the meat. Still others have suggested placing meat in a closed container containing liquid marinade and subjecting the container to a vacuum on the basis that a vacuum will cause the marinade to penetrate more deeply into the meat. U.S. Pat. No. 2,525,581 includes the concept of placing meat in a liquid within a vessel and altering the pressure applied to the meat during a marination process for more effective penetration of the marinade.\nWhile forcing a marinade into meat achieves the goal of improving, or at least changing, the taste of the meat after it has been cooked, marination alone does not dramatically effect how tender the meat will be when cooked.\nWhen the strength of connective tissue is great and of substantial quantity, or the muscle tissue is strong and is not broken down in the cooking process, meat can be \"tough\", that is, difficult to chew. Accordingly, it has long been an objective of processors of meat to find ways of enhancing the tenderness and the flavor thereof.\nExamples of the work done by others to improve the tenderness, texture or flavor of meat are exemplified by the following United States patents:\n(a) U.S. Pat. No. 4,818,550 teaches applying a vacuum to draw marination into the meat but does not teach rapid variation of a high vacuum to achieve tenderization; PA1 (b) U.S. Pat. No. 2,043,392 relates to tenderizing meat but in this method a dilute aqueous solution of a proteolytic enzyme is pumped through the vascular system of the meat; PA1 (c) U.S. Pat. No. 2,509,299 relates to a method of preserving food in which the food is placed is a syrup composition including a bleach and subjecting the composition while in the syrup to an elevated fluid pressure and freezing the food while subjected to the elevated pressure; PA1 (d) U.S. Pat. No. 2,881,080 is for tenderizing meat by applying an aqueous enzyme solution and then subjecting the meat to sound wave energy; PA1 (e) U.S. Pat. No. 5,200,223 is a process for treating meat by subjecting the meat to thin jets of liquid to penetrate the meat in which the jets are applied in short bursts; and PA1 (f) U.S. Pat. No. 5,328,403 subjects meat to shock waves generated by an explosive charge to tenderize the meat.\nWhile all of these previously disclosed methods of enhancing the tenderness and flavor of meat have merit to one degree or another, nevertheless, none adequately address the problem of economically and expeditiously tenderizing meat by breaking down excessive connective and muscle tissue. Particularly none have taught the concepts to effectively tenderize meat with the opportunity at the same time to marinate the meat to produce a superior flavor when the meat is cooked.\nFor further background information concerning the treatment of meat, see the following additional United States patents:\n______________________________________ PATENT NO. INVENTOR TITLE ______________________________________ 2,525,581 Bierman Apparatus For Treating Food Material 2,742,367 Apparatus and Method For Curing Meats 3,370,959 Method of Flavoring Poultry Meat 3,663,233 Method of Tenderizing, Curing and Cooking A Meat Product 4,036,122 Apparatus For Treating Meat, More Particularly Ham Meat 5,057,332 Apparatus and Process For Marinating Foodstuffs ______________________________________\nOf the known prior art references, U.S. Pat. No. 2,525,581 appears to be the most closely related to the present invention."}
{"text": "The present invention relates to an electrically erasable and programable random access read only memory (hereinafter called as an EEPROM).\nA conventional EEPROM is disclosed in U.S. Pat. No. 4,901,320 (Literature I) or Japanese laid open publication 64-59693 (Literature II). Literature I relates to an error correction principle and its method, whereas Literature II relates to an EEPROM operable with low voltage and low current. An EEPROM circuit with high liability has recently been demanded. In other words, it has been desired to achieve an EEPROM circuit having relatively longer life time and less failure which can operate with low voltage and low current and can detect its failure."}
{"text": "This invention relates to variable area exhaust nozzles for gas turbine engines and, more particularly, to sealing means for nozzle flaps of turbojet engines.\nThe exhaust nozzle of a gas turbine engine, such as a turbojet or turbofan engine, has as its purpose a transformation of the pressure and thermal energy of the combustion discharge into velocity, with the forward thrust of the engine being directly proportional to the increase in velocity of the gas from the entrance of the engine to the exit plane of the nozzle. In high performance engines and, in particular, in engines having some sort of thrust augmentation such as an afterburner, it has been found desirable to cause a variation of nozzle flow area to maintain high engine performance under a wide range of operating conditions. For example, it is desirable to provide a larger nozzle flow area during a take-off mode of operation than during a cruise mode. In addition to the function of maintaining the exhaust gas temperature within allowable limits, variable area exhaust nozzles may be employed to bring about noise, thrust and fuel economy benefits. One means for varying the nozzle flow area is by the so-called iris mechanism wherein a plurality of concentrically disposed movable members or flaps are pivotably supported about the nozzle axis. One of the problems associated with such an arrangement is the need to maintain effective sealing between the flaps as the flaps are adjusted to vary the nozzle flow area. Therefore, it is desirable to provide an exhaust nozzle whose area can be flexibly varied between minimum and maximum positions while maintaining a circumferentially continuous aerodynamic structure throughout the entire range.\nEarly method of locating seals with respect to exhaust nozzle flaps relied entirely on a combination of bolts and spectacles wherein, when the nozzle was in the closed position the seals were relatively free to move in the circumferential direction, and when the flaps moved toward the open position, the position of the seals was still not positively enough controlled so as to maintain circumferential sealing integrity around the entire nozzle periphery. Some of the problems encountered were those of dimensional stack-up, limited seal overlap within the circumferential envelope, and misalignment due to nozzle sag on or near the horizontal plane. These problems caused nozzle leakage and seal \"blow-out\", thereby resulting in decreased nozzle efficiency.\nRecent methods of effecting positive placement of seals within exhaust nozzles employ a combination of linkage pairs interconnecting the flaps to the seal wherein an axial track is located on the seal for the purpose of varying the effective lengths of the links. Such an arrangement has been recognized as being somewhat complex and requiring an excess number of moving parts which are susceptible to wear and malfunction.\nAccordingly, a primary object of the present invention is to provide an improved seal arrangement for a jet engine variable exhaust nozzle flap.\nAnother object of this invention is the provision in a variable exhaust nozzle for the maintaining of circumferential sealing integrity throughout the range of nozzle areas.\nYet another object of the present invention is the provision for maintaining a variable area exhaust nozzle seal in a centered relationship between adjacent flaps during all modes of nozzle operation.\nThese objects and other features and advantages become more readily apparent upon reference to the following description when taken in conjunction with the appended drawings."}
{"text": "This invention relates to the field of hermetic ceramic substrates for microelectronic applications. More particularly, this invention relates to protective surface layers for providing hermeticity and corrosion protection.\nIn the microelectronics environment, there is a need for high density, high strength packaging to provide interconnection between semiconductor devices and connection from the devices to the electrical power supply. The electrical properties which are desirable include a highly conductive medium in a highly insulative carrier medium having a low dielectric constant. Thermally, the package must withstand not only the operating environment but also the thermal excursions encountered during the processing and fabrication of the part.\nMechanically, it is preferable to have a substrate package which can withstand chip and pin joining stresses and stresses related to interconnecting with the next level of packaging.\nThe packaging should also be hermetic to prevent degradation of any of the desired properties due to adverse environmental effects. In particular, corrosion of any of the metallurgies used in the package is a real concern if not adequately protected from ambient moisture and deleterious ionic contaminants.\nBoss et al. U.S. Pat. No. 4,880,684, the disclosure of which is incorporated by reference herein, discloses the present day state of the art for the bottom surface of a ceramic substrate. Thus, there is a capture pad, a polymeric sealing and stress relief layer followed by an I/O (conventional nomenclature meaning \"input/output\") bonding pad. Thereafter, the I/O pin is brazed to the I/O bonding pad.\nIt has recently been discovered by the present inventors, however, that there is a concern with the structure disclosed by Boss et al. The I/O bonding pad is typically made from a plurality of layers of metallic material. The present inventors have found the presence of corrosion at this I/O bonding pad during reliability testing which has led to premature failure of the I/O pad structure. While not wishing to be held to a particular theory, it is believed by the present inventors that the corrosion of the I/O bonding pad occurs by a galvanic corrosion mechanism wherein the less noble metals become anodic to the more noble metals in the presence of a suitable electrolyte such as water. At the edges of the I/O bonding pad, the different metals of the pad are exposed together to the environment which thus allow these deleterious corrosion cells to form.\nIt has thus been proposed by the present inventors to apply a protective coating to the bottom surface of the ceramic substrate after the pins have been attached to protect the edges of the I/O bonding pads.\nVarious solutions have been proposed to protect electronic substrates from the effects of corrosion.\nBakos et al. U.S. Pat. No. 4,048,356, the disclosure of which is incorporated by reference herein, discloses in general the application of a hermetic topsealant for the active areas of an integrated circuit device.\nDarrow et al. U.S. Pat. No. 4,233,620, the disclosure of which is incorporated by reference herein, discloses the application of an epoxy material to the surface of a substrate having electrically conductive pins protruding therefrom and to the sides of the substrate so as to hermetically seal the backside of the substrate containing a chip in a cap. Such arrangements wherein epoxy is applied are typically not reworkable, a disadvantage for many ceramic substrates such as those contemplated by the present inventors wherein reworkability is a necessity. Additionally, epoxies are not suitable for the substrates contemplated by the present inventors because epoxies are brittle which could lead to their cracking. Manifestly, such cracking destroys any advantage epoxies may have as a barrier material. The poor thermal stability of epoxies is yet another disadvantage.\nDalencon U.S. Pat. No. 4,360,559, the disclosure of which is incorporated by reference herein, discloses the application of a protective varnish to a printed circuit card, including the weld area of the pins.\nHarris U.S. Pat. No. 4,427,715, the disclosure of which is incorporated by reference herein, discloses an inorganic passivation layer such as polysilicon glass or vapox overlapping the edge of the bonding pad that is applied prior to the bonding operation. The patent is directed to a structure for TAB bonding pads on semiconductor chips that is aimed at preventing the cracking of the chip during the thermo-mechanical bonding process.\nClark et al. U.S. Pat. No. 4,592,944, the disclosure of which is incorporated by reference herein, discloses a polymeric coating over the top surface of substrates for corrosion protection, insulation, etc. The corrosion protection is directed at the top surface thin film circuitry and solder connections.\nNotwithstanding the numerous solutions and disclosures proposed by those familiar with electronic substrates, there still remains a very real need to solve the problem discovered by the present inventors, namely, the corrosion of the I/O bonding pad due to corrosive effects.\nAccordingly, it is an object of the present invention to solve the problem of I/O bonding pad corrosion by providing a protective layer over the I/O bonding pads to insulate them from the deleterious effects of corrosion.\nThis and other objects of the invention will become apparent to those skilled in the art after referring to the following description considered in conjunction with the accompanying drawings."}
{"text": "The present invention relates to the driving of a floppy disk and, more particularly, to the driving of a tunnel erase type floppy disk.\nA conventional tunnel erase type floppy disk comprises a read/write (R/W) head and a pair of erasing heads positioned to the rear of the read/write head. The pair of erasing heads are operated for erasing both sides of data recorded by the read/write head, so that a non-data record area is provided between the tracks to assure compatibility.\nNormally, because the record tracks are arranged to be coaxial, the outer portions of the record disk move more slowly than the inner portions of the disk. Each of the tracks contain an identification (ID) field for storing address information, a data field for storing data, and a gap. If while recording data into the floppy disk the read/write head and the pair of erasing heads are simultaneously operated and the track is improperly aligned, the pair of erasing heads may erroneously erase the data in the ID field, and may not conduct a tunnel erasure of the data field. Therefore, conventionally, the pair of erasing heads are operated as being delayed at a predetermined time after the read/write head is energized.\nBecause the outer tracks of the disk move more slowly than the inner tracks of the disk, a common predetermined delay time is not suitable for all the tracks.\nDue to the usage of a common delay time, it may be difficult to delay the ON condition of each of the pair of erase heads at a satisfactory time with respect to all of the tracks."}
{"text": "The present invention relates to the pumping art, and, more specifically, concerns a new and improved construction of hose pump containing movable pressure elements.\nIn the case of hose pumps, a hose is rhythmically compressed in sections for the purpose of conveying a fluid medium. To this end the hose pump is either equipped with one or a number of squeeze rolls, arranged at a rotatable body, and a hose arranged at a circular arcuate-shaped surface, or with a number of adjacently arranged punches which are moved to-and-fro by cams and a flat or planar counter element. Conveying of the fluid medium is accomplished either by the rotating squeeze rolls or by successive contact of the punches against the counter element, so that the hose is continuously pressed in a peristaltic movement.\nExperience has shown that the hose of such type hose pumps are exposed to large loads and subjected to pronounced mechanical wear. This wear is attributable to the continuously identical deformation of the hose, i.e. the hose is always compressed together at the same region, and the additional tensional load of the hose by the rotating squeeze rolls, or by the additional loading by virtue of the edges of the punches."}
{"text": "Database management systems (DBMS) provide centralized structured information services to client programs. The data within a DBMS is defined, stored, and retrieved according to certain basic information storage patterns known as data models. For example, a DBMS using a hierarchal data model would relate data entities within it through a tree-like structure with each entity having a single “parent” entity (or no parent, for the “root” entity) and any number of “child” entities. This hierarchal data model is used in the popular XML data format and for XML databases.\nThe most common DBMS data model for sophisticated computerized information is the relational data model, which is used by relational database management system (RDBMS). Relational data can be understood as a set of two-dimensional data tables with rows representing data entities and columns representing entity properties. Rows typically have a key column that identifies each row or entity uniquely. These keys are most often a positive integer incremented with each subsequently inserted row. Because table rows may be identified solely through their keys, a row in table A can “relate” to a row in table B by including a column with a key from table B.\nBy using these row identifiers for relations between tables, all data for a row, other than the row identifier itself, does not need to be duplicated and may be stored only in its canonical table. For example, a table of customers may contain street address information. When referring to this customer, such as in a list of sales orders, only the customer row identifier is required, as described before; the customer address does not need to be duplicated within the customer portion of the sales order. One of the major advantages of these normalized intra-table relations is that the customer address may be modified by only changing its information in the customer table as there are no duplicate inclusions of the customer address (i.e. in the sales orders) to maintain.\nRow identifiers are typically automatically generated by the RDBMS as a simple incrementing positive number such as 1, 2, 3, and so on. As these identifier sequence generators are distinct from the tables themselves, if two databases share the same types of tables and columns, the row identifiers may be different between the two databases despite having identical content. These differences can arise from the order, number, and deletion history of the table rows. For example, if three customers were added in opposite order in two databases, likely zero, but no more than one row would have identical row identifiers across the databases.\nWhile an RDBMS provides benefits of reduced storage size, improved data consistency, and, for some operations, higher query performance, its data model is extremely dependent upon correct and consistent use of the row identifier keys. Since data duplication is reduced or even eliminated, if the correct row is not related through the row identifier, there is no recourse to recover the correct row by examination of included duplicate data. For example, a sales order would contain only the customer row identifier and not additional customer columns such as their name or street address which would allow cross-referencing. Furthermore, mistaken row identifiers have a high likelihood of not merely referring to missing rows in another table, but may relate extant, but incorrect rows.\nA DBMS typically runs on a single server in a single location, although more esoteric varieties are sometimes transparently spread among locations. As a RDBMS contains two principle parts, (1) the data tables describing and defining the data structure, and (2) the data rows which conform to the model so described, a server may likewise be seen as both the defining container and the content so contained. While interaction between completely distinct databases (i.e. both the tables and the data are distinct) is possible, there are a number of situations where multiple RDBMS servers are operational with the same table definitions, though the data they hold may be the same or different. In these same-definition scenarios, there may be a need to transport data between the database servers.\nOne such scenario is representing a multitude of software installations through a multitude of database servers. For example, stores in multiple locations may all have the same sales-tracking software installed, but each have their own database server to track their local sales data. In this case, one may want to move partial data, such as a customer and their related records such as order history, from one installation to the server of another installation. Another scenario involves duplicating the same data in several different database servers, perhaps located in different regions, to allow continued availability and security. A third scenario is the use of data within one database server as the template for new data insertion in another server. A variation of this scenario is the use of such a template within the original database server.\nFrom the aforementioned particulars about relational databases, there emerge four major difficulties in accomplishing the above transportation of tabularly identical data between and within an RDBMS:    (1) The transported data row identifier cannot be safely used in the new database due to collisions with existing identifiers and missing relations.    (2) The transported row may be the required relation or require relation with additional rows that should be transported as well. For example, transporting a product row may require supplier (relation to the product) or support history (relation from the product) rows.    (3) As the previous row identifier cannot be safely used, new row identifiers must be allocated as part of transportation.    (4) Given these new, different row identifiers, all necessary related data transported with the primary row must be re-related using the new row identifiers.\nBecause of these difficulties, transportation between relational databases is typically restricted to three scenarios, all limited in practice. These three limited scenarios are:    (1) A database may be completely duplicated or “cloned” into a second database. This is a destructive process in which all existing data in the second database must be completely erased. This is the typical technique by which a database is backed up and restored.    (2) A custom program may be written to collect the specific rows necessary and then insert them into the second database. This technique allows for partial data transport, but is entirely specific to the section of data, database, and intended application as it requires a unique program for each section of data transported and must be synchronously maintained to match any changes to the table definition or intention.    (3) In certain cases the data may be exported out to an intermediate format, such as CSV (Comma Separated Value) format, and then reimported. As the intermediate format typically lacks the full expression of relations and dependencies, this scenario is of limited use except for very simple data transportation such as for rows that are neither related to nor from which are related.\nAs these scenarios all demonstrate considerable restrictions of specificity, application, or scope, there presently exists no general solution for transportation of relational data, in particular for partial interrelated subsets of a RDBMS."}
{"text": "Recursive sequence recombination entails performing iterative cycles of recombination and screening or selection to “evolve” individual genes, whole plasmids or viruses, multigene clusters, or even whole genomes (Stemmer, Bio/Technology 13:549-553 (1995)). Such techniques do not require the extensive analysis and computation required by conventional methods for polypeptide engineering. Recursive sequence recombination allows the recombination of large numbers of mutations in a minimum number of selection cycles, in contrast to traditional, pairwise recombination events.\nThus, recursive sequence recombination (RSR) techniques provide particular advantages in that they provide recombination between mutations in any or all of these, thereby providing a very fast way of exploring the manner in which different combinations of mutations can affect a desired result.\nIn some instances, however, structural and/or functional information is available which, although not required for recursive sequence recombination, provides opportunities for modification of the technique. In other instances, selection and/or screening of a large number of recombinants can be costly or time-consuming. A further problem can be the manipulation of large nucleic acid molecules. The instant invention addresses these issues and others."}
{"text": "1. Technical Field\nThis invention relates to a pill dispenser. More particularly, this invention relates to a compartmentalized pill dispenser that is small enough to be carried with the user at all times, seals the pills against moisture and other degrading elements, protects the pills against vibration and is easy to use by an impaired person in an emergency situation.\n2. Background Art\nThere are many instances where a person is required, due to illness, to carry medicine, such as pills and capsules, with them to be taken at prescribed times or upon the indication of need. An example of such a situation would be a person who has angina pectoris. In the case of an angina attack, an emergency arises due to deficient oxygenation of the heart muscles. This emergency requires an immediate dosage of nitroglycerin to be orally administered by the person affected. Since time is of the essence a person must, as of necessity, carry the pills on their person at all times and administer them immediately when an emergency situation arises. In such an emergency situation, a lapse of even a few minutes without administration of the medication could be harmful or even fatal.\nIt is conventional practice that persons who are required to carry pills with them at all times to keep these medications in their original containers. One drawback of this practice is that many such containers are bulky and are not easily carried in one's pockets. Another drawback that can result in dire consequences is that a person must first remove the cap of the container, remove the cotton stuffing and then single out a one small pill for ingestion. Often this problem is made more difficult if the person requiring the medicine has trembling or shaking hands or the container has a child-proof cap. Additionally, such an emergency situation may occur at night or while the victim is driving, or poor eyesight may make it particularly difficult for him or her to single out the tablet or pill and take it before their condition worsens. If the pills are spilled the user may not be able to find them in time to prevent harm.\nIt is also known for person's required to carry pills with them to employ a small container. For example, nitroglycerin tablets are often carried in a small tube with an inside diameter just larger than the pills themselves. Usually five to seven pills are stacked one on top of another in this tube. Unfortunately, these containers are known to fail such as when the cap which seals the tube becomes cross-threaded and stuck thereby making it difficult to remove the cap. Sometimes the pills are crushed by the cap if too many are loaded within the dispenser. If space is left to avoid this problem vibration of the pills increases. Vibration has been known to powderize the nitroglycerin pills to the point where they become jammed within the tube and cannot be extracted. Additionally, when the pills become powderized the dosage is then uncertain. Another problem with these tube dispenser is that like the larger bottles, the pills stacked within the tube can be easily spilled; this is especially true when the user is having some kind of attack. The dispenser is also difficult to load, especially for someone with poor sight, because it is often difficult to ascertain the number of pills within the container.\nAdditionally, many pill dispensers of the prior art do not protect the pills or tablets against moisture, light or other degrading effects. This is important as some medications are detrimentally affected by these environmental factors. For example, nitroglycerin degrades substantially when exposed to moisture or light and the pills are easily pulverized when exposed to vibration.\nOne attempt at resolving the access and spilling problems involved a pen-shaped dispenser, having an elongated hand-held core with a plurality of cavities adapted to contain pills or tablets. A sheath, open at both ends, surrounds the core which is slideable from a position covering all of the cavities to another position uncovering the cavities. A detent mechanism is employed between the sheath and the core for interrupting the sliding movement after each of the cavities is uncovered. For spatially orienting the tablet dispenser, and without the need to look at it, distinct physical touch indicia are fixed to the dispenser in lengthwise alignment with the tablet cavities. Although this pill dispenser resolved some of the problems associated with carrying pills on the person at all times, other persist. For instance, although the core and the sheath of the container are in close proximity to each other, there is not a water-tight seal between the core and the sheath of the container. This allows air and water to enter the dispenser and degrade the pills or tablets over time. Additionally, although the container is smaller than some of the prior art, its size is still substantial when considering its user would have to carry it with them constantly. Its length approximates that of a pen--on the order of six inches in length. In fact, the container purposely had an extension and pen clip such that the container could be carried in one's shirt pocket. Since women, and many men, usually do not place pens in any shirt pocket they might have, this would render the dispenser somewhat impractical for a person to carry with them.\nTo prevent the common pill dispenser problems, a better pill dispenser than is provided by the current technology is desirable. Ideally, this pill dispenser should be easy to use even when the user is in an impaired state. It should also be light-weight and small enough to take with one at all times. It should be able to dispense individual doses of medicine without the user having to touch them. Additionally, it should protect the pills from the effects of moisture, sunlight, air, and vibration. It is also necessary that the pill dispenser be made of materials that do not react with the pills thereby contaminating them."}
{"text": "1. Field of the Invention\nThe present invention relates to an endoscope in which two or more kinds of adapters are individually detachably attachable to a distal end portion located at a distal end in an inserting direction of an insertion section.\n2. Description of the Related Art\nAs is well known, endoscopes are widely used in a medical field and an industrial field. With an endoscope used in the medical field, it is possible to, by inserting an elongated insertion section into a body cavity serving as a subject, observe an organ in the body cavity and perform various kinds of treatment using a treatment instrument inserted into a treatment instrument insert-through channel according to necessity.\nWith an endoscope used in the industrial field, it is possible to, by inserting an elongated insertion section into a jet engine, a pipe in a factory, or the like serving as an object, perform observation, various kinds of treatment, and the like of scratches, corrosion, and the like in the object.\nA configuration is well known in which, in a distal end portion located on a distal end side in an inserting direction of an insertion section of an endoscope (hereinafter simply referred to as distal end side), an image pickup unit including an objective lens unit and an image pickup device such as a CCD is provided and an illumination unit or the like that illuminates an inside of an object is provided.\nA configuration is also well known in which the objective lens unit and the illumination unit are provided in a known optical adapter (hereinafter simply referred to as adapter) detachably attachable to the distal end portion of the insertion section.\nAs the adapter, a direct-view adapter for observing a front in the inserting direction of the insertion section (hereinafter simply referred to as front) and a side-view adapter for observing a side different from the inserting direction are well known. The respective adapters are properly used according to observation targets and uses.\nFurther, a plurality of adapters having different visual field angles is present in each of the direct-view adapter and the side-view adapter. These adapters are also properly used according to observation targets and uses. Therefore, a plurality of kinds of adapters is individually detachably attachable to the distal end portion.\nJapanese Patent Application Laid-Open Publication No. 2004-313241 discloses a configuration in which, in a plurality of adapters detachably attachable to a distal end portion, identification resistors having different resistance values for each of the adapters are provided and, when an adapter is attached to the distal end portion, an identification resistor comes into contact with a terminal for resistor identification provided in the distal end portion and energizes the identification resistor, whereby a CPU electrically connected to the terminal via a lead wire in a control unit, to which an endoscope is connected, reads a resistance value and automatically detects a type of the adapter attached to the distal end portion from the resistance value.\nJapanese Patent Application Laid-Open Publication No. 2004-313241 also discloses a configuration in which, since it is difficult to reduce the adapter or the distal end portion in diameter if a light source such as a light emitting element is provided in the adapter or the distal end portion of an insertion section, the light source is provided in the control unit, a light guide that guides illumination light irradiated from the light source to a distal end face located at a distal end in an inserting direction of the distal end portion (hereinafter simply referred to as distal end) is inserted through the endoscope, the illumination light emitted from a distal end of the light guide is made incident in the adapter, and an illumination optical system that supplies the illumination light into an object is provided.\nIncidentally, in general, the attachment of the adapter to the distal end portion including the configurations disclosed in Japanese Patent Application Laid-Open Publication No. 2004-313241 is performed by coating, with an inner circumference of a proximal end side in the inserting direction of the adapter (hereinafter simply referred to as proximal end side), an outer circumference of a mounting and detaching section located on a distal end side of the distal end portion and screwing, for example, with a male screw formed in the outer circumference of the mounting and detaching section, a female screw formed in an inner circumference of a stop ring extending from a proximal end in the inserting direction of the adapter (hereinafter simply referred to as proximal end) backward in the inserting direction (hereinafter simply referred to as backward).\nIn the distal end portion including the mounting and detaching section, built-in components such as the lens unit in the image pickup unit, the terminal functioning as a member for adapter identification explained above, the lead wire extending from the terminal, the light guide explained above, and various channels are provided."}
{"text": "In rolling bearings, the sliding tracks of the rolling elements must be hardened. In bearings with a double row of rolling elements, in particular with a double row of balls used in vehicle hub bearing units, induction hardening of the rolling tracks of the outer ring of the bearing is a particularly delicate operation. In fact, the presence of sharp edges between the two adjacent tracks implies risk of over-heating and cracking; on the other hand, if the heating power is decreased to avoid these risks, there is the risk of not achieving a hardening profile that is suitable in depth, extension or surface hardness of the hardened area.\nMoreover, the induction heating devices, or more simply, inductors currently used, are bulky, difficult to construct and/or maintain and sometimes with a not entirely satisfactory reliability."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an injection apparatus; particularly an injection apparatus maintaining the nozzle and the injection gate at respective desired temperatures.\n2. Background of the Invention\nIt has long been known that the temperature of a melt material is important to successful injection. This is particularly true when the melt material has a high melt temperature. For example, polyethylene terephthalate (\"PET\") is typically injected above 500.degree. F. A drop in the temperature of the melt material prior to reaching the injection cavity would lower the melt material temperature below that required for proper melt material flow causing less than ideal flow characteristics. These flow characteristics can cause deformed or defectively molded parts; particularly when injecting multilayer parts comprising very thin layers. Therefore, it is desirable to maintain the nozzle temperature at or above the temperature required to assure proper melt material flow as the melt material leaves the nozzle.\nIt is also known to maintain an injection cavity at a temperature relatively low compared to the temperature of the melt material to facilitate quick cooling of the melt material upon reaching the cavity. The colder the cavity temperature at the time the melt material is injected, the faster the melt material will solidify and allow removal of the solidified part from the cavity. Therefore, a relatively lower cavity temperature will decrease the overall cycle time for injection molding a part. Moreover, it is known that if the injection gate temperature exceeds the desired temperature of the melt material, `stringing` of the melt material will occur in the nozzle and gate area as the injected part is removed from the cavity after injection is complete. These `strings` either break off with the injected part and interfere with further processing of the part (e.g. blowmolding) or stay in the gate or cavity and cause a physical or aesthetic defect in subsequently injected parts.\nFor these reasons, it has been found desirable to prevent excessive heat transfer from the injection nozzle to the injection cavity. The melt material can thus be maintained at its appropriate temperature in both the nozzle and the cavity. Prior injection apparatuses were often designed to space a nozzle tip from an associated injection cavity during injection to leave a gap therebetween. It was thought that this gap would act as a thermal break between the nozzle and the cavity and allow the nozzle to operate at high temperatures while maintaining a relatively cool cavity. Unfortunately, the thermal break of this configuration could not be maintained at efficient cycle times. During the injection process, melt material would deviate from the injection path and flow into the gap between the nozzle and the cavity. The thermal break thus became a thermal bridge.\nOther attempts to insulate an injection nozzle from a cavity have involved the use of nozzle inserts. For example, U.S. Pat. No. 4,279,588 issued to Gellert and entitled \"Hot Tip Seal\" disclosed a seal (12) located between the nozzle and the injection gate to limit heat transfer therebetween. The seal (12) of Gellert resided substantially within the nozzle and extended outward therefrom to contact the cavity. Similarly, U.S. Pat. No. 4,521,179 issued to Gellert and entitled \"Injection Molding Core Ring Gate System\" disclosed a nozzle seal (76). The seal (76) of Gellert also resided substantially within the nozzle and extended outward therefrom to contact the cavity.\nIt has been found that movement of the various parts within an injection apparatus will result from thermal expansion as portions of the apparatus are heated from ambient temperature to the temperature necessary to inject a melt material. Different injection apparatuses accommodate this thermal expansion in different ways. It has been found that the thermal expansion of some injection apparatuses results in movement of the nozzle both along the longitudinal axis thereof and perpendicular to that longitudinal axis. In other words, it has been found that the nozzles of some apparatuses will elongate and shift laterally as the apparatus is heated. Seals that attached to the nozzle, such as those of the Gellert patents discussed above, break or deform due to this lateral nozzle movement. Such seals are therefore inapplicable to apparatuses experiencing this lateral nozzle movement.\nIt has also been found that many seals cannot withstand the high temperatures and pressures associated with injection; especially when the high temperatures are maintained for long periods of time. Many prior inserts degraded after prolonged exposure to high temperatures resulting in rupture or deformation of the inserts which allowed melt material to leak into the area between the nozzle and the cavity causing in a thermal bridge.\nIt has also been known to supply a cooling means to a cavity to remove the heat transferred from the nozzle or melt material to the cavity. Cooling ducts circulating coolants such as glycol were typically employed. However, the distance between the part void and the injection gate has heretofore limited the proximity of the cooling ducts to the injection gate."}
{"text": "1. Field of the Invention\nThe present invention relates to a fire extinguisher, and more particularly, to a fire extinguisher which is simple in structure and easy to operate.\n2. Description of the Related Art\nA fire extinguisher having a cylinder filled with fire-extinguishing gas is constructed such that when a seal plate of the cylinder is manually or automatically broken at the outbreak of a fire, the fire-extinguishing gas is ejected from the cylinder to put out the fire. Many of conventional fire extinguishers contain component parts that are used in common in manual fire extinguishers and automatic fire extinguishers. Fire extinguishers of this type are, however, disadvantageous in that they are complicated in structure, require a large number of parts, and are costly. Further, some fire extinguishers require complicated handling or involve a difficulty in accurately aiming at the origin of a fire, which hinders fire-fighting work or lessens the fire-extinguishing capability.\nAlso known is a fire extinguisher having a replaceable cylinder. In this type of fire extinguisher, the component parts thereof must be reset in their initial state when the cylinder is replaced; this often makes it difficult to smoothly load a new cylinder. If the cylinder cannot be smoothly replaced, then it is difficult to perform fire-fighting work using a plurality of cylinders successively replaced."}
{"text": "Not applicable\nThe present invention is directed to an OFDM/DMT digital communications system. More particularly, the present invention is directed to an apparatus and method for synchronizing the clocks used in a transmitter and receiver of an OFDM/DMT digital communications system. The present invention is particularly applicable in multipoint OFDM/DMT digital communications systems.\nMulti-point communications systems having a primary site that is coupled for communication with a plurality of secondary sites are known. One such communications system type is a cable telephony system. Cable telephony systems transmit and receive telephone call communications over the same cable transmission media as used to receive cable television signals and other cable services.\nOne cable telephony system currently deployed and in commercial use is the Cablespan 2300 system available from Tellabs, Inc. The Cablespan 2300 system uses a head end unit that includes a primary transmitter and primary receiver disposed at a primary site. The head end unit transmits and receives telephony data to and from a plurality of remote service units that are located at respective secondary sites. This communication scheme uses TDM QPSK modulation for the data communications and can accommodate approximately thirty phone calls within the 1.9 MHz bandwidth typically allocated for such communications.\nAs the number of cable telephony subscribers increases over time, the increased use will strain the limited bandwidth allocated to the cable telephony system. Generally stated, there are two potential solutions to this bandwidth allocation problem that may be used separately or in conjunction with one another. First, the bandwidth allocated to cable telephony communications may be increased. Second, the available bandwidth may be used more efficiently. It is often impractical to increase the bandwidth allocated to the cable telephony system given the competition between services for the total bandwidth available to the cable service provider. Therefore, it is preferable to use the allocated bandwidth in a more efficient manner. One way in which the assigned bandwidth may be used more efficiently is to use a modulation scheme that is capable of transmitting more information within a given bandwidth than the TDM QPSK modulation scheme presently employed.\nThe present inventors have recognized that OFDM/DMT modulation schemes may provide such an increase in transmitted information for a given bandwidth. Such systems, however, present a number of technical problems. One such problem is the determination of how one or more remote receivers are to synchronize their internal clocks and timing systems with the internal clock and timing system of a primary transmitter at a central site. A remote receiver must first synchronize its internal clock and timing system with the clock used by the primary transmitter to synthesize the transmitted signal before the remote receiver can properly demodulate the data that it receives. A further problem occurs in multipoint communication systems in which there are plural groups of remote transmitters that transmit data to a central transceiver. Each group of transmitters often has its transmissions frequency multiplexed with transmissions from other groups before being demultiplexed for receipt by the central transceiver. The resulting multiplexing/demultiplexing operations introduce frequency offsets for which compensation must be made if the receiver of the central transceiver is to properly extract the correct data from the signals that is receives. The present inventors have recognized the need for such upstream and downstream clock synchronization and have disclosed solutions to these problems.\nIn a communications system comprising a transmission medium, symbols are generated in a predetermined number of bins using at least a first timing signal. A reference signal also is generated by using said first timing signal. The symbols and said reference signal are transmitted across a transmission medium carried by a carrier signal. The generating the transmitting are preferably achieved with a transmitter. At a first point along the transmission medium, the carrier signal with a second signal are frequency multiplexed, preferably by a frequency multiplexer. The carrier signal and the second signal are transmitted to a second point on the transmission medium displaced from the first point. The carrier signal is frequency demultiplexed after the carrier signal and second signal have reached the second point, preferably by a demultiplexer. The demultiplexed carrier signal is frequency demodulated in response to the reference signal, preferably by a first demodulator. The symbols and the reference signal are demodulated in response to the frequency demodulated carrier signal, preferably by a second demodulator.\nOther features and advantages of the present invention will become apparent upon review of the following detailed description and accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates generally to integrated circuit semiconductor memory devices and associated methods. In particular, the present invention relates to a column select line (CSL) control for which the same signal controls the enable timing and the disable timing signals for synchronous random access memory devices.\n2. Description of Related Art\nSpeed improvements in semiconductor memory devices, such as Dynamic RAMs and Static RAMs, have historically come from process and photolithography advances. More recent memory speed improvements, however, have resulted mainly from making changes to the base architecture. An example of fast RAM architecture is the synchronous architecture. One key advancement of the synchronous memories is their ability to synchronously burst data at a high-speed data rate. Additionally, in a system with a synchronous RAM, since data, addresses and control signals are latched into the memory in synchronism with the system's clock signal, the system's processor is able to perform other tasks freely until data is available after a known number of clock cycles. This architecture provides substantial advantages in memory operating performance.\nIn a typical semiconductor memory device, in order to write/read data into/from a specific memory cell in a memory device, the specific memory cell should be designated by a row address and a column address. When the specific memory cell is designated in a read/write operation, a charge distribution operation is performed with respect to data read out from the designated memory cell to a bit line, and the readout data is amplified by a sense amplifier. The amplified data is transmitted to an input/output line through an I/O gate circuit, and then is output from the memory chip via associated output circuits. The read operation of one-bit data stored in the specific memory cell is completed by the above-described process. The column decoder turns on the selected I/O gate by receiving and decoding the column address.\nTo simplify the complexity of the decoding operation in highly integrated memories, a column pre-decoder is typically provided to pre-decode the column address prior to the main decoding operation therefor. This column decoding scheme has been adopted in most high density memory devices.\nFIG. 1 is a block diagram illustrating a conventional exemplary synchronous memory device. Referring to FIG. 1, an array 100 of memory cells is provided to store data. Word lines WL0-WLm and bit lines BL0-BLn coupled with the cells run along the rows and columns of the memory cell array 100, respectively. In the vicinity of the cell array 100, a row decoder 120 is provided for selectively driving the word lines WL0-WLm, and an input/output (I/O) gate circuit 140 for supporting the selective transmission of data from the bit lines BL0-BLn to a data I/O buffer 280, and vice versa. The I/O gate circuit 140 is controlled by column select lines CSL0-CSLn. Externally applied address signals A0-Ax including both column and row address signals are fed to an address buffer 160. The column address signals CA0-CAi among the address inputs are applied to a column pre-decoder 180.\nA clock buffer 230 is suppled with an external clock signal XCLK and provides an internal PCLK synchronized with the external clock signal XCLK. A CSL enable control circuit 240 generates a CSL enable control clock signal PCSLE by the logical combination of the internal clock signal PCLK and a column address setting signal PYE from a timing control logic (not seen). The column pre-decoder 180 pre-decodes the column address signals CA0-CAi and generates pre-decoded address signals DCA0-DCAj.\nThe column pre-decoder 180 outputs the DCA0-DCAj signals under the control of the PCSLE signal from the CSL enable control circuit 240. Main decoding operation of the column address signals are then carried out by a column main-decoder 200. This decoder 200 generates decoded signals DCAB0-DCABk by decoding the DCA0-DCAj. The DCAB0-DCABk signals are provided to a column driver 220 which drives the column select lines CSL0-CSLn selectively in response to the DCAB0-DCABk signals. A CSL disable control circuit 260 generates a CSL disable control clock signal PCSLD by the logical combination of the internal clock signal PCLK and a normally logic-high signal PVCCH. The column driver 220 is disabled by the PCSLD signal from the CSL disable control circuit 260, and hence stops driving the column select lines CSL0-CSLn.\nFIGS. 2A and 2B illustrate the constructions of the CSL enable and disable control circuits 240 and 260, respectively, in detail. Referring first to FIG. 2A, the CSL enable control circuit 240 includes a delay circuit formed by inverters IV1-IV4 (\"first\" delay circuit), a NAND gate G1, and an inverter IV5. The internal clock signal PCLK is provided to the delay circuit. The NAND gate G1 has one input applied with the delayed signal of the clock signal PCLK and the other input applied with the column address setting signal PYE. The output signal of the NAND gate G1 is output through the inverter IV5 as the CSL enable control clock signal PCSLE.\nReferring to FIG. 2B, the CSL disable control circuit 260 includes a delay circuit formed by inverters IV6-IV8 (\"second\" delay circuit) and a NAND gate G2. The second delay circuit is also fed with the clock signal PCLK. This delay circuit has a smaller delay time than the first delay circuit. The output of the second delay circuit is supplied to one input of the NAND gate G2. The normally logic-high signals PVCCH is provide to the other input of the NAND gate G2. This gate G2 outputs the CSL disable control clock signal PCSLD.\nFIG. 3 shows the detailed configuration of a unit circuit of the column pre-decoder 180. As shown in FIG. 3, the unit pre-decoder circuit 180' includes inverters IV31-IV49 and NAND gates G34-G49. The unit column pre-decoder circuit 180' is provided with three column address signals CA0-CA2 from the address buffer 160, and generates eight pre-decoded column address signals DCA0-DCA7. The CSL enable control clock signal PCSLE is commonly applied to the first inputs of the NAND gates G42-G49. The second inputs of the NAND gates G42-G49 are provided with the substantial pre-decoded column address signals, i.e., the output signals of the inverters IV34-IV41, respectively. When the PCSLE signal becomes high, the output signals of the inverter IV34-IV41 can be propagated to the inverters IV42-IV49 via the NAND gates G42-G49, respectively, and they are output as the pre-decoded column address signals DCA0-DCA7. The PCSLE signal should go high only after the completion of the pre-decoding operation with the inverters IV31-IV41 and the NAND gates G34-G41 in order to prevent pre-decoding errors.\nFIG. 4 illustrates the detailed construction of unit circuits of the column main-decoder 200 and column driver 220, respectively. Referring to FIG. 4, the unit column main-decoder circuit 200' includes NAND gates G50-G57, and inverters IV50-IV57 corresponding to the NAND gates G50-G57, respectively. Each of the NAND gates G50-G57 has one input provided with a corresponding pre-decoded column address signal DCAy (where, y=0, 1, . . . , or 7) and the other input with a gate control signal GCS from a timing control logic (not shown). Each output signal of the NAND gates G50-G57 is provided as a finally decoded signal DCABy (where, y=0, 1, . . . , or 7) through the corresponding inverter IV50, IV51, . . . , or IV57.\nThe unit column driver circuit 220' includes inverters IV60-IV67, cascode inverters (sometimes called \"dual gate inverters\") 40-47, and inverting latches 60-67. Each of the cascode inverters 40-47 consists of two PMOS transistors (e.g., MP40a and MP40b) and one NMOS transistor (e.g., MN40), and each inverting latch (e.g., 60) is formed of two cross-coupled inverters (IV60a and IV60b). For each cascode inverter (e.g., 40), three transistors (MP40a, MP40b and MN40) have their source-drain paths coupled in series between a boosted supply voltage terminal VEXT and a ground voltage terminal GND. Each decoded column address signal (e.g., DCAB0) is applied to the gates of the corresponding pull-up and pull-down transistors (MP40a and MN40). The PCSLD signal from the CSL disable control circuit 260 is commonly fed to the gates of the switching transistors MP40b-MP47b of the respective cascode inverters 40-47 via the inverters IV60-IV67. Each inverting latch (e.g., 60) is coupled to the drain junction of the corresponding switching and pull-down transistors (MP60a and MN60).\nThe pre-decoding operation begins with the CSL disable control clock signal PCSLD of a high level. A pre-decoded column address signal of a high level (e.g., DCA0) can be transferred to the gate of the PMOS transistor MP40a as a decoded column address signal DCABO only when the gate control signal GCS remains at a high level. Namely, the GCS signal determines whether to propagate the DCA0-DCA7 signals through the unit column main-decoder circuit 200' or not. The decoded signal DCAB0 goes high when both DCA0 and GCS signals are high, so PMOS pull-up transistor MP40a turns off and NMOS pull-down transistor MN40 on. The high-level DCAB signal is latched by the inverting latch 60, so that a corresponding column select line CSL0 is driven high. After the PCSLD signal has gone low, the GCS signal also goes low. Accordingly, the pull-up transistor MP40a turns on and the pull-down transistor MN40 off, but the column select line CSL0 still remains high owing to the inverting latch 60. In this situation, when the PCSLD signal goes high again, the switching transistor MP40b turns on, so the CSL0 line is driven low.\nAs described above, the column select lines CSL0-CSLn are selectively activated by the column pre-decoder 180, but deactivated by separately controlling the column driver 220.\nFIG. 5 is a timing diagram illustrating read/write operations of the conventional synchronous memory device of FIG. 1. With reference to FIG. 5, after a column address strobe signal CAS is activated low, in clock cycle T0, the CSL disable control clock signal PCSLD goes high in synchronism with the external clock signal XCLK (or the internal clock signal PCLK). After a predetermined time (i.e., Tm1) has elapsed, during which the first column address signals CA#0 (i.e., CA0-CAi) have reached the column pre-decoder 180, the CSL enable control clock signal PCSLE goes high in response to the activation of the column address setting signal PYE (see FIG. 2A). Of course, the PCSLE signal is also synchronized with the clock signal XCLK (or PCLK). A unit column pre-decoder circuit 180' pre-decodes the column address signals CA#0 (CA0-CA2) and generates the pre-decoded column address signals DCA#0 (DCA0-DCA7) of which only one is active and the others inactive. Here, assuming DCA0 signal is activated high, then a corresponding column select line CSL0 will be driven high by a unit column driver circuit 220'.\nIn the next clock cycle T1, the PCSLD signal becomes high before the low-to-high transition of the PCSLE signal, so that the line CSL0 is deactivated. Next, after the second column address signals CA#1 (CA0-CA2) had reached the unit column pre-decoder circuit 180' and the time Tm1 has elapsed, when the CSL enable control clock signal PCSLE goes high again in response to the activation of the column address setting signal PYE. The unit column pre-decoder circuit 180' generates the second decoded column address signals DCA#1 (DCA0-DCA7). Here, assuming DCA1 signal is activated high, then a corresponding line CSL1 will be driven high by a unit column driver circuit 220'.\nThe other column select lines (such as CSL2 and CSL3) also will become activated and deactivated during next clock cycles (T2 and T3) in response to the other column address signals (such as CA#2 and CA#3) in the same manner as the above-mentioned.\nIn the above conventional memory device, the CSL enable control clock PCSLE should not go active until valid column address signals arrive at the pre-decoder 180 during each clock cycle Tc (where, c=0, 1, 2, . . . ). However, in the event the PCSLE signal goes high during a clock cycle (e.g., T1) before the valid column address signals CA#1 arrive at the column pre-decoder 180, owing to an insufficient delay time of Tm1, then the invalid column address signal CA#0 for the previous clock cycle T0 may be pre-decoded again by the corresponding column pre-decoder circuit 180' (refer to FIG. 3). Hence, the invalid decoded signal DCAB0 will be latched by the corresponding inverting latch 60 via the cascode inverter 40 (refer to FIG. 4). This leads to the activation of the column select line CSL0. Thereafter, when the valid decoded signal DCAB1 is activated by decoding the valid column address signals CA#1 and latched by a corresponding inverting latch 41 in cycle T1, the column select line CSL1 corresponding to the valid column address signals CA#1 also becomes active along with the invalid CSL0 line, causing an erroneous read/write operation. For the above reason, it is essential to ensure sufficient delay time Tm1 in the conventional memory device. This limits the memory access speed improvements.\nIn addition, according to the conventional memory device structure, a significant area penalty may result from the large reiterative layout area of the unit column driver circuit 220'.\nFurthermore, since the pull-up and switching transistors MP40a-MP47a and MP40b-MP47b within the respective cascode inverters 40-47 provide current leakage paths together with the inverters IV60b-IV67b of the inverting latches 60-67 during power-up, the conventional device has large power-up current dissipation."}
{"text": "1. Field of the Invention\nThe present invention relates generally to producing maleic anhydride and more particularly to shaped catalyst structures containing catalytic material comprised of mixed oxides of vanadium and phosphorus which are useful for the production of maleic anhydride.\n2. Background of the Invention\nVarious shapes for tablets of vanadium and phosphorus catalysts used in the production of maleic anhydride from butane have been described in the literature. For example, U.S. Pat. Nos. 4,283,307 and 5,168,090 describe catalyst tablets of various shapes that are used for the production of maleic anhydride.\nOne disadvantage of current tablet shapes is that they have relatively low side crush strengths. Side crush strength measures the weight or force needed to break a tablet. The side crush strength is important in catalyst manufacture since it directly impacts the durability of the catalyst during handling, shipment, and installation into a commercial reactor. With all other factors being equal, a catalyst tablet with a higher side crush strength is preferred to one with a lower crush strength.\nAnother disadvantage of current tablet shapes is that they have relatively high percent attrition. The percent attrition is a measurement (as detailed later) of the amount of catalyst lost to disintegration after subjected to a certain amount of wear and tear. The percent attrition is important in catalyst manufacture since it also directly impacts the durability of the catalyst during handling, shipment, and installation into a commercial reactor. With all other factors being equal, a catalyst tablet with a lower percent attrition is preferred to one with a higher percent attrition.\nCatalyst tablets with lower crush strength and higher attrition tend to produce more broken catalyst tablets and more catalyst fines during handling, shipment, and installation into a commercial reactor. These broken catalyst pellets and catalyst fines create increased pressure drop in the commercial reactor tubes during normal operation and are thus undesirable."}
{"text": "The present invention relates generally to a feedforward amplifier for use mainly in the high-frequency band and, more particularly, to a feedforward amplifier of increased power efficiency.\nFIG. 1 illustrates in block form a basic configuration of a feedforward amplifier. The feedforward amplifier comprises a distortion detector 13 made up of an amplifying transfer path 11 of a main amplifier (hereinafter referred to as a main amplifier path) and a first linear signal transfer path (hereinafter referred to as a first linear path) 12 and a distortion canceller or suppressor 16 made up of a second linear signal transfer path (hereinafter referred to as a second linear path) 14 and a distortion injection path 15. The main amplifier path 11 is formed by a series connection of a variable attenuator 21, a variable phase shifter 22 and a main amplifier 23. The first linear path 12 is formed by a series connection of a delay line 24 and a phase inverter 25. The second linear path 12 is formed by a delay line 26. The distortion injection path 15 is formed by a series connection of a variable attenuator 27, a variable phase shifter 28 and a auxiliary amplifier 29. The input into the feedforward amplifier is divided by a power divider 31 to the main amplifier path 11 containing the main amplifier 23 and the first linear path 12. A power combiner/divider 32 generates the sum of and the difference between signals from the main amplifier path 11 and the first linear path 12, and provides them to the second linear path 14 and the distortion injection path 15, respectively. The output from the feedforward amplifier is produced by a power combiner 33 which combines outputs from the second linear path 14 and the distortion injection path 15.\nAccordingly, the feedforward amplifier detects a distortion component (a difference component) generated by the main amplifier 23 in the distortion detector 15 and, in the distortion canceller 16, regulates the phase and amplitude of the distortion component and injects it into the output signal from the main amplifier 23 provided via the second linear path 14, thereby canceling nonlinear distortion generated by the main amplifier 23. In general, the amount of nonlinear distortion cancelled by the feedforward amplifier depends on the regulation of the variable attenuator 21 and the variable phase shifter 22 of the distortion detector 13 and the variable attenuator 27, the variable phase shifter 28 and the auxiliary amplifier 29 of the distortion canceller 16. The accuracy of such regulation is disclosed in Japanese Patent Laid-Open Gazette No. 1-198809 entitled \"Automatic Regulation Circuit for Feedforward Amplifier.\" For example, phase and amplitude deviations for achieving a distortion compression of 30 dB or more are .+-.2.degree. and .+-.0.3 dB, respectively. It can be said, therefore, that strict conditions are imposed on the balance and completeness of regulation of the transmission characteristics of the distortion detector 13 and the distortion canceller 16.\nThe feedforward amplifier compensates for the nonlinear distortion generated by the main amplifier 23. Accordingly, because of its circuit configuration the feedforward amplifier cannot theoretically compensate for nonlinear distortion by the auxiliary amplifier 29. Furthermore, because of the stringent condition on the balance of each of the above-mentioned two paths, linearity of input vs. output power performance is required of the auxiliary amplifier 29 of the conventional feedforward amplifier. To increase the linearity of the amplification circuit using a semiconductor amplifying device, it is customary in the prior art to use a Class A bias for the condition of operation to thereby make the saturation output voltage sufficiently larger than the peak voltage of the signal to be amplified.\nIn recent years there has been a growing demand for small, light and inexpensive radio equipment of lower-power consumption. The same is true of radio equipment using the feedforward amplifier. To reduce the power consumption of the feedforward amplifier, it is essential to increase the power efficiency of both of the main and auxiliary amplifiers. With the increased power efficiency of the both amplifiers, it is possible to reduce their cooling bodies and so on and hence achieve downsizing of radio equipment.\nThe power efficiency of the main amplifier can be enhanced as by a push-pull circuit which operates under Class B bias condition. The conventional feedforward amplifier is capable of compensating for the nonlinear distortion that occurs in the main amplifier. On the other hand, to increase the power efficiency of the auxiliary amplifier 29 inserted in the distortion injection path 15 of the feedforward amplifier, it is necessary, in general, that the semiconductor amplifying device forming the auxiliary amplifier be caused to operate under a Class B or C bias condition. The nonlinear distortion that results from such bias condition cannot theoretically be compensated for by the feedforward amplifier as referred to above. For these reasons, enhancement of the power efficiency of the auxiliary amplifier 29 gives rise to the problem of impairing the distortion compensating ability of the feedforward amplifier.\nThe power-supply efficiency of the feedforward amplifier can be expressed by the ratio between the output power of the feedforward amplifier and the power supplied thereto. For example, according to Toshio Nojima and Shouichi Narahashi, \"Ultra-Low, Multi-Frequency Common Amplifier for Mobile Communications, --Self-Adjusting Feedforward Amplifier (SAFF-A)\" Technical Report of Institute of Electronics, Information and Communication Engineers of Japan, RCS90-4, 1989, in the case where the saturation output power of the main amplifier is 100 W, the saturation output power of the auxiliary amplifier is 1/8 that of the main amplifier, GaAs-MESFETs are used as semiconductor amplifying devices for the main amplifier and the auxiliary amplifier, the drain voltage and current of the MESFET of the main amplifier are 12 V and 20 A, the drain voltage and current of the MESFET of the auxiliary amplifier are 12 V and 5 A and the both amplifiers are operated in a 1.5 GHz band under the Class A bias condition, the power supply to the feedforward amplifier is 300 W. If mean output power back off is set at 8 dB and the loss of the main amplifier output signal in the distortion canceller is ignored, then the feedforward amplifier output is approximately 15 W. Accordingly, the power-supply efficiency of the feedforward amplifier is in the range of between 15/300 and about 5%. Even in the case where a highly efficient amplification circuit such as a Class B push pull amplifier is used as the main amplifier and a Class A amplifier as the auxiliary amplifier, the power-supply efficiency is around 10% or less at the highest.\nTo achieve high power efficiency of a high-output power amplifier, there is described in W. H. Doherty, \"A new high efficiency power amplifier for modulated waves,\" Pro. IRE, Vol. 24, No. 9, pp.1163-1182, September 1936 a method which employs a plurality of amplifiers of different saturated output voltages. This method is generally known under the name of Doherty method, which has been implemented, for example, in a medium-wave transmitting power amplifier for broadcasting stations. In the Doherty method, a saturation amplifier and a linear amplifier are connected in parallel. The saturation amplifier amplifies signals of mean power or thereabout and the linear amplifier amplifies signals of peak power. The Doherty method achieves high power efficiency amplification by means of the saturation amplifier, but because of the circuit configuration used, signals that ought to be input into the linear amplifier are also applied to the saturation amplifier--this raises the problem of nonlinear distortion. Moreover, the circuit configuration by this method is incapable of compensating for the nonlinear distortion caused by the saturation amplifier.\nOne method that has been proposed to further increase the powersupply efficiency of the feedforward amplifier is to set the output backoff (a difference between the saturated output power and the output power at the operating point) of the main amplifier at a small value. As is well-known in the art, however, the nonlinear distortion that can be compensated for by the feedforward amplifier is limited only to incompleteness of the linear region of the input/output power characteristic of the main amplifier. That is, in general, the feedforward amplifier cannot compensate for nonlinear distortion as by clipping in the saturation region of the input characteristic."}
{"text": "The invention relates to an apparatus for transmitting a movement between two components, which, as a function of a first, translational movement of a first coupling element to which a first component is coupled, produces a translational movement of a second coupling element to which a second component is coupled.\nSuch an apparatus is known from DE 195 24 636 C1, for example as an angled drive for a high-voltage power breaker. This document describes the capability to drive three switch poles of a three-pole high-voltage power breaker by means of a common drive via a drive linkage, with each switch pole being connected to the drive by means of a crank which, firstly, changes the direction of the drive movement, and secondly produces a step-up ratio.\nIn this document, the different initial positions of the individual cranks result in a time offset in the operation of the individual switch poles.\nEach individual angled crank comprises a two-armed lever whose lever arms are at right angles to one another, with the ends of the lever arms forming coupling elements to which, firstly, a vertical switching rod and, secondly, the horizontal drive rod can be coupled as components. Any translational movement of the free end of the first lever arm, which is connected to the drive rod, results in a translational movement of the free end of the second lever arm, which is connected to the switching rod.\nThe described design results in the switch poles switching in a staggered manner one after the other in time.\nThe present invention is based on the object of providing an apparatus of the type mentioned initially, which allows a variable step-up ratio using simple means.\nAccording to the invention, the object is achieved in that the first and the second coupling element are connected to one another, such that they can move, by means of a first rigid element, and in that the first coupling element is connected in a hinged manner to a first lever which is mounted at a first, fixed-position bearing point such that it can pivot, and in that the second coupling element is connected in a hinged manner to a second lever which is mounted at a second fixed-position bearing point such that it can pivot, and in that the first and the second bearing point are spaced apart from one another, and are located on opposite sides of the movement path of the first rigid element.\nWhen a movement is being transmitted, the apparatus according to the invention results in the first and the second component, which are coupled to one another by the apparatus, each being moved at essentially the same speed in a first phase of the movement, in the movement travel of the second component lagging behind that of the first component in a second phase, and in the second component being accelerated in a third phase. When a movement is being transmitted in the opposite movement sense, the second component is first of all accelerated with respect to the first component and, after this, the second component is retarded in a corresponding manner.\nAdvantageous refinements of the inventions can be found in the dependent claims 2 to 7.\nThe fixed-position bearing points are advantageously each located in the plane of the movement paths of the rigid elements and of the coupling elements.\nThe coupling elements are each arranged in the end regions of the rigid elements and may be in the form of connecting bolts or of some other mechanical component, to which the various other parts such as the lever and rigid elements can be coupled. The coupling elements may also be in the form of a linkage.\nThe invention also relates to a high-voltage power breaker having a number of switch poles, at least two of which can be driven by a common drive, with at least one of the common driveable switch poles being connected to the drive by means of an apparatus as claimed in one of claims 1 to 5.\nA high-voltage power breaker has an interrupter unit in each switch pole, in which interrupter unit two switching contacts are brought into conductive contact with one another, or are disconnected from one another. For this purpose, at least one switching contact can be driven by means of the drive.\nSuch a high-voltage power breaker can easily be adjusted, for use in a polyphase high-voltage system having an apparatus according to the invention, in such a manner that the individual switch poles have a specific time offset, during switching on, with regard to the making of the first conductive contact between the switching contacts, for example with a lead or lag of 3 milliseconds between one another in each case. In this way, the various switch poles can be operated at a respectively advantageous time, with regard to the respective phase of the polyphase alternating current to be switched on.\nIt is also possible to use a common drive to operate two of the switch poles simultaneously for capacitive switching, while the third switch pole is operated with a time offset of 5 milliseconds after them. The third switch pole is then connected to the common drive by means of an apparatus according to the invention. For such offset switching operations, the individual switch poles are relieved as well as possible of the load involved in controlled operation in synchronism with a mains system."}
{"text": "A typical carrying strap, for example, a carrying strap for a rifle, golf bag, is constructed with a strap of leather or canvas as the primary load-bearing member and a wider and thicker layer of foam or other suitable padding placed between the strap and the wearer's shoulder. A typical carrying strap sometimes tends to slip or roll from the shoulder. Anti-slip carrying straps are available on the market. The anti-slip carrying straps available on the market tend to prevent slip in one direction, e.g., lateral, but not in other directions, e.g., longitudinal. Accordingly, there is a continuing need for a pad that prevents the carrying strap from slipping off the shoulder."}
{"text": "Radio frequency identification (RFID) is a data collection technology that enables the storing and remote retrieval of data utilizing devices referred to as RFID tags, or transponders. An RFID transponder may comprise a silicon integrated circuit, or chip, and an antenna that enables the RFID transponder to receive and respond to radio frequency (RF) queries from an RFID transceiver, or reader. The RFID transponder may comprise memory, for example a random access memory (RAM) or an electrically erasable programmable read only memory (EEPROM), which enables storage of data. The data may comprise an electronic product code (EPC) that may be utilized to locate an item to which the RFID transponder is attached. For example, libraries may attach RFID transponders to books to enable the tracking of books that are checked out to library patrons. RFID transponders may be integrated into “smart cards” and/or into other electronic devices such as mobile phones. The RFID transponders may enable storage of account information that enables the holder of the smart card to purchase goods and services. The transponder may, for example, store a current balance that indicates a monetary value of goods and services that may be purchased with the transponder. The transponder holder may purchase goods and services by holding the transponder in the proximity of an RFID reader that retrieves account information from the smart card.\nDesigners of RFID transponders are challenged with the tasks of increasing the range at which the transponder can communicate with the reader and improving the energy efficiency due to the limited power available to the transponder, all while keeping the transponder extremely small and inexpensive.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "1. Field of the Invention\nThis invention relates to apparatus for pushing stacks or packs, particularly packs of tubular film sections, from a support, preferably from a pallet, on which said stacks or packs have been deposited in a plurality of layers, comprising a pusher, which is movably mounted in a frame and moved by a drive.\n2. Description of the Prior Art\nIn the manufacture of sacks, a tube-making apparatus is operated to eject tubular film sections at a rate which is much higher than the rate at which said tubular films can be processed by the succeeding end-forming machine. Even if a tube-making machine is succeeded by a plurality of end-forming machines it may be necessary to effect an intermediate storage of the tubular film sections which have been ejected by the tube-making machine and have been stacked to form packs. Such intermediate storage is usually effected by means of a storage belt conveyor. As such belt conveyors also have only a limited capacity it may be necessary to deposit packs of tubular film sections on pallets.\nThe packs of tubular film sections which have been deposited in a plurality of layers on pallets must be returned to the production process and the tubular film sections must be singled and then supplied to the end-forming machine.\nSingling may be effected, e.g., by so-called rotary feeders, which comprise stack magazines to which the packs of tubular film sections are supplied. Such rotary feeder which serves to single the tubular film sections and comprises a conveyor for supplying the packs of tubular film sections and for introducing them into the stack magazine of the rotary feeder is known from Published German Application No. 30 46 280. The apparatus known from Published German Application No. 30 46 280 is also provided with a pusher, by which the packs of tubular film sections supplied by a storage belt conveyor are pushed onto the receiving end of the conveyor which forwards them to the stack magazine.\nIn apparatus of the kind disclosed in German Patent Specification No. 31 05 154 an intermediate storage of packs of tubular film sections is effected on carrying grates, which are vertically spaced apart to form a plurality of tiers, and said packs are pushed from such grates by means of a pusher onto a conveyor in case of need; the latter conveyor serves to return the packs of tubular film sections into the production process.\nThe known apparatus can only be used to push packs of tubular film sections from a conveyor or a storage location onto another storage location or a conveyor. But when storage space is no longer available for an intermediate storage of juxtaposed packs of tubular films, then it will be necessary to stack the packs of tubular film sections in a plurality of vertically spaced apart planes, for instance, on pallets. Thus far, the removal of packs of tubular film sections stacked, e.g., on pallets and the deposition of said packs on a conveyor for returning said packs to the production process had to be effected by hand."}
{"text": "1. Field of the Invention\nThe invention relates to a method of making measurements on and between existing objects located underwater. More particularly, the invention relates to controlling the drift of an inertial navigation system using a range and bearing image capture system such as a laser camera or similar sensors, to make measurements without the system being in contact with an underwater object or the seabed. The spatial relationship between the inertial navigation system and the underwater object is achieved by range and bearing measurements.\n2. Description of the Prior Art\nAny motion of an object in space is composed at the microscopic level of two basic motions: micro-rotations and micro-linear motions. An inertial navigation system (hereinafter “INS”) measures these motions using gyroscopes (to sense angular motions) and accelerometers (to sense linear motions). The INS may measure rotations using the gyroscope and linear motions using the accelerometers at the high-speed microscopic level (typically 500 times per second). Even given this high sample rate, the INS is subject to error. Because of the imperfection of the accelerometers and gyroscopes, TNS navigation—when not controlled—is subject to a continuous drift that reaches about 0.8 miles per hour in a more accurate INS and several miles per hour in one that is less accurate.\nThis drift is especially problematic when the INS is used in underwater metrology i.e., the science of measuring accurately the position and orientation of objects and physical structures underwater (often on the seabed). The most common metrology consists in measuring the distance between the two flanges of a pipeline system that will receive the ends of a jumper or spool i.e. the hard, prefabricated connecting pipes that join a pipeline system (typically ending at a manifold) and a wellhead. The accurate measurement of distance between receptacles of the pipeline and the wellhead, to which each end of the jumper is to be attached, is crucial to ensure accurate manufacturing of the jumper. To that end, the exact distance between the pipeline receptacle and the wellhead receptacle, the difference in depth between them and the exact 3D orientation (heading, pitch and roll) of each receptacle must be accurately measured.\nTo correct for drift in underwater metrology, the INS navigation is sometimes corrected using an “aiding by zero velocity update”, or zupt method. The zupt is a period without motion where the INS is parked on some structure on the seabed: usually either one of the 2 receptacles that need to be measured. When the INS is parked for at least 15 seconds, the INS recalibrates and corrects any navigation drift from the knowledge that when there is no motion the readings of the accelerometers and gyroscopes must be caused by noise (pure erroneous measurements). The draw back of such a method is that it requires that the INS carried by a robot (such as a Remotely Operated Vehicle or “ROV”) or other vehicles or a diver, touches the flange or receptacle (actually: the INS rests on the receptacle for at least 15 seconds), potentially causing damage. Even worse, in order to measure accurately the position and orientation of those flanges, a mating system (stabbing guide) is often specially made and installed on or near the two flanges to perform these traditional metrologies which is cumbersome.\nAccordingly, embodiments of the present method and device, control the drift of an INS used in underwater inertial metrology without touching any seabed structure, and therefore without the requirement for stabbing guides. Thereby the invention allows for “touchless” metrology."}
{"text": "1. Field of the Invention\nThe present invention relates to a garden protector and more particularly pertains to a device which may be used to scare off animals from an area under protection without damage to such animal.\n2. Description of the Prior Art\nThe use of automatically activated water sprinklers is known in the prior art. More specifically, such devices heretofore devised and utilized for the purpose of automatically turning the sprinkler on are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements. Such devices are utilized primarily for watering a lawn or garden at preset times although U.S. Pat. No. 4,253,606 differs in that it turns a sprinkler system on when the humidity in the soil reaches a predetermined low point. Other devices operating off timers are shown in U.S. Pat. Nos. 5,074,468; 3,865,138; 3,669,352; and 3,599,669.\nIn this respect, the device according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in so doing provides an apparatus primarily developed for the purpose of turning a sprinkler on only when the need to scare off an intruder in a protected area exists. Such a device eliminates danger of injury to a child or animal such as may occur when other types of devices such as electric fences are used for area protection.\nTherefore, it can be appreciated that there exists a continuing need for new and improved sprinklers which can be used as protective devices. In this regard, the present invention substantially fulfills this need."}
{"text": "1. Technical Field\nThe present disclosure relates to a surgical stapling apparatus. More particularly, the present disclosure relates to an endoscopic surgical stapling apparatus that includes a mechanism for articulating a tool assembly.\n2. Background of Related Art\nSurgical fastening devices, such as surgical stapling apparatuses, are well known in the art. These devices typically include two elongated jaw members to capture or clamp tissue. One jaw member carries a staple cartridge that houses a plurality of staples arranged in at least two lateral rows while the other jaw member has an anvil that defines a surface for forming the staple legs as the staples are driven from the staple cartridge. For the most part, the stapling operation is effected by cam members that travel longitudinally through the staples cartridge, with the cam members actuating upon staple pushers to sequentially eject the staples from the staple cartridge. A knife may travel between the staple rows to longitudinally cut or open the stapled tissue between the rows of staples. Examples of this kind of instrument are disclosed in U.S. Pat. Nos. 3,079,606 and 3,490,675.\nU.S. Pat. No. 3,499,591 also discloses a stapling apparatus with a pair row of staples on each side of the incision. The patent discloses a surgical stapler that has an upper jaw member and a lower jaw member. The upper jaw member includes a disposable cartridge while the lower jaw member contains a disposable anvil having a series of staple-closing depressions aligned with a plurality of staple-carrying grooves of the cartridge. In operation, a cam member translates through an elongated guide path between the two sets of staple-carrying grooves. Staple drive members are located within the grooves and are positioned in such a manner so as to effect ejection of the staples from the staple cartridge. Other examples of such stapling instruments are disclosed in U.S. Pat. Nos. 4,429,695 and 5,065,929.\nEach of the instruments described hereinabove is designed for use in conventional surgical procedures wherein surgeons have direct manual access to the operative site. Endoscopic or laparoscopic procedures, however, are performed through a small incision or through a narrow cannula inserted through a small entrance wounds in the skin. To address the specific needs of endoscopic or laparoscopic surgical procedures, surgical stapling devices have been developed and are disclosed in, for example, U.S. Pat. Nos. 5,040,715; 5,312,023; 5,318,221; 5,326,013; 5,332,142; and 6,241,139.\nTyco Healthcare Group, LP, the assignee of the present application, has successfully manufactured and marketed endoscopic stapling instruments, such as the Multifire ENDO GIA™ 30, Multifire ENDO GIA™ 60 and Multifire ENDO TA™ 30 instruments, for a number of years. An articulating endoscopic stapler known as the ENDO GIA™ Universal stapler has also been marketed. Typically, these devices include an articulation mechanism having an articulation lever operatively engaged with a cam member having a stepped camming channel. Current stapling instruments allow manual operation of the articulation lever. These instruments have provided significant clinical benefits. Nonetheless, improvements to these instruments are still possible.\nIt would be desirable to provide an improved articulation mechanism for a surgical stapling apparatus. It would also be desirable to provide an articulation mechanism capable of being operated either manually or electromechanically.\nAccordingly, it is an object of this disclosure to provide an improved articulation mechanism that articulates the tool assembly of a surgical stapling apparatus by electromechanical means. It also the object of this disclosure to provide an improved articulation mechanism that will allow a user to articulate with minimal effort and with only one hand the tool assembly of a surgical stapling instrument.\nAnother object of the disclosure is to provide an articulation mechanism capable of very accurate and fine position adjustment at a speed that can be easily controlled by a variable speed switch.\nStill another object of the disclosure is to provide an articulation mechanism with a manual override for additional reliability and safety."}
{"text": "Semiconductor memory devices are typically classified into volatile memory devices and non-volatile memory devices. Volatile memory devices are subdivided into dynamic random-access memories (DRAMs) and static random access memories (SRAMs). Non-volatile memory types include erasable programmable read-only memories (EPROMS) and electrically erasable programmable read-only memories (EEPROMs). EEPROMs are increasingly used in system programming that requires continuous update or auxiliary memory devices. Particularly, flash EEPROMs are advantageous as mass storage devices because their integration density is high compared with conventional EEPROMs.\nFrequently, it would be convenient to be able to mix integrated circuit device types, such as EEPROMs and other memory devices, and bipolar integrated circuits, such as NPN transistors, onto a single integrated circuit chip. However, due to the inherently low breakdown voltage (approximately 10 volts) of typical wells used in BiCMOS technology and the need for a high programming voltage of an EEPROM memory device (approximately 14 volts), there has been no simple and economical way to integrate these two device types into a single integrated circuit. Previously, the problem has been avoided in the art by using additional masks to create high-voltage wells for EEPROMS."}
{"text": "Electro-luminance displays have been developed for a wide variety of devices, such as cell phones. In particular, active-matrix organic light emitting diode (AMOLED) displays with amorphous silicon (a-Si), poly-silicon, organic, or other driving backplane have become more attractive clue to advantages, such as feasible flexible displays, its low cost fabrication, high resolution, and a wide viewing angle.\nAn AMOLED display includes an array of rows and columns of pixels, each having an organic light emitting diode (OLED) and backplane electronics arranged in the array of rows and columns. Since the OLED is a current driven device, the pixel circuit of the AMOLED should be capable of providing an accurate and constant drive current\nThere is a need to provide a method and system that is capable of providing constant brightness with high accuracy and reducing the effect of the aging of the pixel circuit and the instability of backplane and a light emitting device."}
{"text": "This invention relates generally to pinball games and, more particularly, relates to an improved ball trough for such games.\nPinball games typically consist of an inclined playfield supporting a plurality of play features such as targets, bumpers, and the like. Disposed on the playfield are a pair of flippers which are used to direct a ball at selected game features for scoring points. Play usually begins with the ball positioned in the shooter lane where the player manipulates a spring loaded plunger to propel the ball onto the playfield. Play continues until the ball drains from the playfield via a drain hole which is normally located behind the flippers to capture balls which have been misplayed. Located in the playfield, out of sight of the player, is a ball trough which extends from the drain hole to the shooter lane for transporting balls therebetween.\nAn example of a prior art ball trough may be found in U.S. Pat. No. 5,358,243 to Eddy et al. which is herewith incorporated by reference in its entirety. The ball trough disclosed therein consists of a generally U-shaped channel that is mounted to the underside of the playfield and extends between the drain hole and the shooter lane. The bottom surface of the channel is inclined relative to the horizontal such that a ball entering the drain hole will roll under the force of gravity toward the shooter lane. When activated, a solenoid plunger will kick the ball adjacent the shooter lane up through an aperature in the playfield and into the shooter lane. A series of optical switches are also mounted on the channel and arranged such that each switch will detect the presence or absence of each ball in the queue. In this embodiment, the kicker must supply enough force to the ball to elevate the ball back to the level of the playfield.\nWhile this embodiment is useful in achieving its purpose, it has been noticed that the procedure used to elevate the ball into the shooter lane generates a vibrational force large enough to shake the entire structure. These vibrational forces have been known to causes some of the optical switches associated with the channel to become displaced. As mentioned, these optical switches are used by the microprocessor to asses playing conditions whereby any failure or misalignment of the optical switches will cause a resulting failure in the game control. Therefore, a need exists for providing an improved ball trough capable of moving the ball into the shooter lane without the jolting associated with currently employed kicker mechanisms.\nAs a result of these exising needs, it is an object of the present invention to provide an improved lifting mechanism whereby balls may be smoothly carried from the ball trough to the shooter lane."}
{"text": "The invention involves a floating oil barrier for the purpose of containing oil or similar materials in water. It is comprised of larger floating partial cells and are made of waterproof materials. Smaller partial cells that are submerged in the water are made of untreated or hydrophilic materials, producing stabilizing effects of the barrier's position, wherein the partial cells are comprised of glass fibers or glass fiber mats and a flexible tube that encases them. They are joined to one another via a traction cable.\nWith larger oil tanker vessels being developed, environmental danger posed by maritime oil spills is also continuously increasing. In recent years, oil disasters have resulted in entire coastlines becoming severely polluted and marine life in the endangered areas being destroyed to such an extent that these large ships are being burdened with greater and greater charges. For this reason, efforts have been made to build multiband tankers, in an attempt to make it more difficult for oil spills to occur. In the end, however, in addition to the associated production problems, such measures cannot prevent oil from escaping if a ship sinks. One method, known in the art from FR-PS 2 646 189, calls for waterproofed mineral fibers to be blown onto the surface of the water. The oil is supposed to settle onto these fibers, so that the bulk can then be picked up as a whole and disposed of. The fibers are waterproofed with silicone oil, and, because they are short, can be hindered from sinking. Contamination and pollution of the sea bed is thus such a problem that mineral fibers of this type are not used. In another method, known in the art from GB-PS 1 235 463, oil on top of the water is to be adsorbed onto inorganic fibers, with the fibers being treated beforehand with a water-repellent material. As a result, the fibers float on top of the oil and cause the oil to settle on them; the oil is then removed by burning off the fibers. The burning off of the oil adsorbed by the fibers is not economically feasible and causes considerable environmental damage. In addition, too little space is available for storage of the fibers, which are simply to be sprayed onto the water. DE-PS 41 40 247 contains the recommendation that waterproofed glass wool materials be joined together to form a barrier, with the glass fibers being inserted into a perforated flexible tube. The oil can enter the tube through the perforations, thus entering the glass wool where the oil is adequately adsorbed and held until the barrier can be removed from the water and disposed of. The disadvantage of this method is that the oil adsorber, or the barrier, floats on the surface of the water. When waves hit the barrier, the oil cannot be kept from splashing back out. Accordingly, DE-OS 43 01 761 or DE-OS 43 25 806 recommends that a multisection oil adsorber be employed, in which one cell is coated with waterproofed material and the others, which lie below this one, are coated with non-waterproofed material. The upper partial cell, coated with waterproofed glass fiber material, floats on top of the water and can absorb oil, while the lower, or several lower, partial cells are submerged in the water, as they more or less attract the water due to their particular formation. They act as a kind of centerboard that effectively prevents the oil from flowing through when waves hit the barrier. Integrated into the joined partial cells is a traction cable, via which the barrier can be maneuvered. In particular, this makes it possible to pull the barrier, which borders a correspondingly large oil slick, back to the ship, thereby continuously skimming off the oil floating on the water or binding this oil using suitable glass fiber mats. Under strong winds, however, when this type of oil barrier is hauled in, it too cannot stop the collected oil from splashing over the barrier and then back into the unprotected marine area, that is, onto the surface of the water. To some extent, eddies also form, which in the end even lead to the oil being pushed under the centerboard, and from there into the oil-free section of the water.\nIt is thus the object of the invention to create an oil barrier which, in strong winds and especially when it is being hauled in, will prevent the oily water from spilling, and thereby will also prevent the oil from passing below the centerboard."}
{"text": "A single row deep groove radial ball bearing according to the present invention is used to support a rotary member such as a pulley and to removably support an offset load.\nAutomotive accessories such an alternator and a compressor are driven to rotate by an engine for driving a vehicle. Due to this, an endless belt is extended between follower pulleys fixed to end portions of rotating shafts of the automotive accessories and a drive pulley fixed to an end portion of a crankshaft of the driving engine, and the rotating shafts are constructed to be driven to rotate based on the circulation of the endless belt.\nFIG. 5 illustrates the construction of a rotational driving portion of a rotating shaft 1 of a compressor constituting an automotive air conditioner. The rotating shaft is rotatably supported by rolling bearings, not shown, within a casing 2. A follower pulley 4 is rotatably supported around the circumference of a supporting tube portion 3 provided on an outer circumference of an end portion of the casing 2 by a single row deep groove radial ball bearing 5. The follower pulley 11 is formed into an annular configuration, on the whole, which has substantially a U-shaped cross section, and a solenoid 6 which is fixed to an end face of the casing 2 is disposed within an internal space of the follower pulley 4. On the other hand, a mounting bracket 7 is fixed to a portion protruding from the casing 2 at an end portion of the rotating shaft 1, and an annular plate 21 of a magnetic material is supported on the circumference of the mounting bracket 7 via a plate spring 8. The annular plate 21 is spaced away from the follower pulley 4, as shown in FIG. 5, when the solenoid 6 is not energized, while when the solenoid 6 is energized, the annular plate 21 is drawn toward the follower pulley 4 so as to be secured thereto, so that a rotational force is free to be transmitted from the follower pulley 4 to the rotating shaft 1.\nWith the rotational supporting device as described above, there may be a case where a transverse central position (a chain line xcex1 in FIG. 5) of the endless belt wound around an outer circumference of the follower pulley 4 is not allowed to coincide with a transverse central position (a chain line xcex2 in FIG. 5) of the single row deep groove radial ball bearing 5. In such a case, a moment load in proportion to a deviating amount (offset amount) xcex4 between the transverse central positions of the two members is applied to the single row deep groove radial ball bearing 5 based on the tension of the endless belt. Then, a central axis of an inner ring 9 and a central axis of an outer ring 10 which constitute the single row deep groove radial ball bearing 5 do not coincide with each other (they are inclined).\nWith a mechanism like one as described above, when the central axes of the inner ring 9 and the outer ring 10 do not coincide with each other, there occurs an unbalanced wear of the endless belt which is wound around the outer circumference of the follower pulley 4, this making it difficult to secure the durability of the endless belt. In addition, the inclination of the central axes also makes it impossible to secure a certain gap between the annular plate 21 and the follower pulley 4, resulting in a possibility that these two members 21, 4 come into friction with each other. In the event that such a friction occurs, abnormal wear and abnormal noise are likely to be generated unfavorably. With a view to preventing the occurrence of these inconveniences, it is considered to reduce an angular gap of the single row deep groove radial ball bearing 5 in order to make it difficult that the central axes of the inner ring 9 and the outer ring 10 discord with each other.\nThen, in order to reduce the angular gap for the aforesaid purpose, the following (1) to (4) procedures will be contrived.\n(1) Radius of curvatures of cross-sectional shapes of an inner ring raceway 11 formed in an outer circumferential surface of the inner ring 9 and an outer ring raceway 12 formed in an inner circumferential surface of the outer ring 10 are made small (they are to be reduced so as to approximate 50% of the outside diameter of balls 13 constituting the single row deep groove radial ball bearing 5).\n(2) As shown in FIG. 6, the raceway surface of one or both of an inner ring raceway 11a in an outer circumferential surface of an inner ring 9a and an outer ring raceway 12a in an inner circumferential surface of an outer ring 10a is formed into a combined surface, and rolling surfaces of balls 11 are brought into contact with both the raceway surfaces at three or four points.\n(3) As shown in FIG. 7, the heights of shoulder portions 14a, 14b existing on transverse (in left and right directions in FIG. 7) sides of the raceway surface of one or both of an inner ring raceway 11b in an outer circumferential surface of an inner ring 9b and an outer ring raceway 12b in an inner circumferential surface of an outer ring 10b are made higher as indicated by a solid line than a general height indicated by a chain line in the same figure.\n(4) As shown in FIG. 8, a plural row radial ball bearing 15 is used in which a plurality of balls 13, 13 are provided between a plurality of inner ring raceways 11c, 11c formed in an outer circumferential surface of an inner ring 9c and between a plurality of outer ring raceways 12c, 12c formed in an inner circumferential surface of an outer ring 10c, respectively.\nThe conventionally known and contrived constructions for reducing the angular gap as described above have the following problems.\nFirst of all, in the case of the construction described under (1), although the angular gap can be reduced, the contact ellipses existing at abutting portions of the rolling surfaces of the respective balls and the inner ring raceway 11 and the outer ring raceway 12 become larger. Then, the contact ellipses dislodge from the inner ring raceway 11 and the outer ring raceway 12 when the central axes of the inner ring 9 and the outer ring 10 are only inclined slightly by virtue of a moment load. In this state, the rolling fatigue life of the rolling surface becomes extremely short. Thus, the construction described under (1) is not desirable as the allowable moment load becomes small. Note that while the configurations of the contact portions can be ellipse no more (resulting in a configuration in which part of the ellipse becomes lost) when the contact portions between the rolling surfaces and the raceway surfaces reach the transverse end edges of the raceway surfaces, for the purpose of description, a state like this will be referred to as xe2x80x9cthe contact ellipse dislodges from the raceway surfacexe2x80x9d in this specification.\nNext, with the construction described under (2), the rolling surfaces of the balls 13 and the inner ring raceway 11a and the outer ring raceway 12a come to contact with each other at a plurality of contact positions, and moreover, in a state in which the engine is driven while the moment load is applied, since the contact positions become asymmetrical relative to the rotating axis of the ball 13, there occur much wear and heat based on slippage at the contact points, which is not desirable.\nNext, with the construction described under (3), since the space between the shoulder portions 14a, 14a on the outer circumferential surface of the inner ring 9b and the shoulder portions 14b, 14b on the inner circumferential surface of the outer ring 10b becomes narrow, the diametrical thickness of a retainer 16 for holding the balls 13 becomes thin. Thus, as the thickness of the retainer 16 becomes thin, since it is difficult to secure the durability of the retainer 16, in consideration of the durability of the retainer 16, the effect of reducing the angular gap using the procedure described under (3) is limited.\nFurthermore, with the construction described under (4), although the effect of reducing the angular gap and securement of the durability of the constituent components becomes compatible at a higher order, the increase in axial dimension cannot be avoided. The rotational supporting portion such as the follower pulley 4 has to be installed within a limited space in many cases, and therefore the increase in axial dimension is not desirable. Moreover, as the axial dimension increases, the production cost of the respective constituent components also increases.\nIt is an object of the present invention to provide a single row deep groove radial ball bearing which can solve the inconveniences as described above.\nSimilarly to single row deep groove radial ball bearings conventionally widely known, any of single row deep groove radial ball bearings according to the invention comprises an inner ring having a deep groove inner ring raceway formed in an outer circumferential surface thereof, an outer ring having a deep groove outer ring raceway formed in an inner circumferential surface thereof and a plurality of balls rotatably provided between the inner ring and the outer ring.\nIn particularly, in the single row deep groove radial ball bearing according to the present invention, at least one of the inner ring raceway and the outer ring raceway has a cross-sectional shape in which a transversely central portion thereof is different in a radius of curvature from transverse end portions between which the transversely central portion is interposed, and also transverse end edges of the portions having different radius of curvatures are made to be smoothly continuous with each other.\nIn the above-mentioned single row deep groove radial ball bearing as set forth in the invention, it is advantageous that the radius of curvature of the cross-sectional shape of at least one (or preferably, both) of the inner ring raceway and the outer ring raceway is made smaller at a transversely central portion and larger at transverse end portions, and transverse end edges of the portions having different radius of curvatures are made to be smoothly continuous with each other.\nAccording to the single row deep groove radial ball bearings constructed as described above in accordance with the invention, the reduction of the angular gap and securement of the allowable moment load can be established at a higher order.\nFirst of all, with the single row deep groove radial ball bearing as set forth in the invention, since the radius of curvature of the cross-sectional shape of the raceway is made larger at the transverse end portions, in the event that the central axis of the inner ring and the central axis of the outer ring are inclined toward each other based on the moment load, and that the contact point between the rolling surface of the ball and the raceway is displaced to the transverse end portion side of the raceway, the contact ellipse existing at the contact point becomes small. Due to this, it is difficult for the contact ellipse to dislodge from the raceway, the allowable moment load can be secured. Moreover, since the radius of curvature of the cross-sectional shape of the transversely central portion of the raceway is made smaller, when compared with a case where the entirety of the raceway constitutes a single curved surface having a large radius of curvature, the rolling surface of the ball and the raceway are allowed to get closer to each other to thereby reduce the angular gap.\nFurthermore, in the above-mentioned single row deep groove radial ball bearing as set forth in the invention, it is advantageous that the radius of curvature of the cross-sectional shape of at least one (or preferably, both) of the inner ring raceway and the outer ring raceway is made larger at a transversely central portion and smaller at transverse end portions, and transverse end edges of the portions having different radius of curvatures are made to be smoothly continuous with each other.\nAccording to the single row deep groove radial ball bearings constructed as described above in accordance with the other aspect of the invention, the reduction of the angular gap and securement of the allowable moment load can be established at a higher order.\nWith the single row deep groove radial ball bearing as set forth in the invention, since the radius of curvature of the cross-sectional shape of the raceway is made smaller at the transverse end portions, when compared with the single curved surface having a large radius of curvature, the rolling surface of the ball and the raceway are allowed to get closer to each other to thereby reduce the angular gap. In this case, too, since the radius of curvature of the cross-sectional shape of the raceway is made larger at the transversely central portion, even in the event that the central axis of the inner ring and the central axis of the outer ring are inclined toward each other based on the moment load, and that the contact point between the rolling surface of the ball and the raceway is displaced to the transverse end portions of the raceway, it is difficult for the contact ellipse to dislodge from the raceway to thereby secure the allowable moment load."}
{"text": "1. Field of the Invention\nThe present invention relates generally to surgical instruments and a methods, and surgical kits for treatment of pelvic floor prolapse conditions.\n2. Background Discussion\nWomen account for more than 11 million of incontinence cases. Moreover, a majority of women with incontinence suffer from stress urinary incontinence (SUI). Women with SUI involuntarily lose urine during normal daily activities and movements, such as laughing, coughing, sneezing and regular exercise. SUI may be caused by a functional defect of the tissue or ligaments connecting the vaginal wall with the pelvic muscles and pubic bone. Common causes include repetitive straining of the pelvic muscles, childbirth, loss of pelvic muscle tone, and estrogen loss.\nNormally, the urethra, when properly supported by strong pelvic floor muscles and healthy connective tissue, maintains a tight seal to prevent involuntary loss of urine. When a woman suffers from the most common form of SUI, however, weakened muscle and pelvic tissues are unable to adequately support the urethra in its correct position. As a result, during normal movements when pressure is exerted on the bladder from the diaphragm, the urethra cannot retain its seal, permitting urine to escape.\nU.S. Pat. No. 5,899,909, which is incorporated herein by reference in its entirety, discloses a surgical instrument comprising a shank having a handle at one end and connecting means at the other end to receive, one at a time, two curved needle-like elements which are connected at one end to one end of a mesh intended to be implanted into the body. In practice, the mesh is passed into the body via the vagina first at one end and then at the other end at one side and the other, respectively, of the urethra to form a loop around the urethra, located between the urethra and vaginal wall. The mesh is extended over the pubis and through the abdominal wall and is tightened. The mesh ends are cut at the abdominal wall, and the mesh is left implanted in the body. This trans-vaginal procedure is exemplified by the TVT product sold by the Gynecare Worldwide, a division of Ethicon Inc., of Somerville, N.J., USA. In this procedure two 5 mm needles pass a PROLENE® mesh trans-vaginally and through the abdomen to create a tension-free support around the mid urethra.\nAn alternate method to treat SUI is the sling procedure. In this procedure a needle or other suture-retrieving device is first inserted through the abdomen, above the pubic bone. The needle is guided behind the pubic bone, through the subrapubic fascia around the urethra, and out of the body through an incision in the anterior vaginal wall. At this point sutures are attached to the needle(s) and pulled up back through the abdominal cavity, where the sutures are fastened to the rectus muscle.\nTechniques for protecting against the puncture of the internal structures during this type of procedure have included laparoscopic procedures. This involves making an incision in the abdomen and inserting a video scope to watch the progress of the needles as they pass through the abdominal cavity. These additional incisions are not optimal for the patient. Also, the needles, which pass through the abdomen, are not designed to capture a mesh but rather a suture, which has been previously attached to the mesh or harvested fascia. These needles are generally in the diameter range of about 0.090 inches to about 0.120 inches. Therefore, the needles do not create a large channel through the fascia. The channel is only wide enough to pass the suture. Accordingly, the sutures do not possess the elongation properties of the PROLENE® mesh and therefore cannot provide the tension-free support of the TVT. Also, attaching a mesh directly to these needles is not optimal because it is very difficult, if at all possible, to pull the mesh through the narrow channel created by the needle.\nAnother common condition suffered by women is prolapse of organs within the pelvic cavity. Prolapse is a condition in which organs, namely the bladder, bowel, and uterus, which are normally supported by the pelvic floor, herniate or protrude into the vagina. This occurs as a result of weakening or damage to the muscles and ligaments making up the pelvic floor support. Childbirth is the most common cause of damage to the pelvic floor, particularly where prolonged labor, large babies and instrumental deliveries were involved. Other factors can include past surgery such as hysterectomy, lack of estrogen due to menopause, and conditions causing chronically raised intra abdominal pressure such as chronic constipation, coughing, heavy lifting, and other physical activity involving impact with the body, such as skydiving. Specific prolapse conditions include cystocele, which is a prolapse of the bladder, rectocele or rectal prolapse, enterocele or intestinal herniation, and prolapse of the vaginal vault such as after a hysterectomy.\nVaginal surgery is the usual method of repair, but abdominal surgery (typically laparoscopc surgery) may also be performed. Traditional pelvic floor repair surgeries, whether abdominal or vaginal, involve lifting the prolapsed organ to restore it back to its correct anatomical position, and subsequently using sutures attached to ligaments and/or muscles to retain the organ in the correct position. Surgeons have also been known to place a layer of mesh below the prolapsed organ, and to subsequently suture corners or sides of the mesh to ligaments or muscles on the sidewalls of the pelvis. The suturing can be done via access through the abdomen or by access through the vaginal incision. Using sutures to support the mesh structure has disadvantages similar to those described above in conjunction with SUI procedures.\nThe device and method disclosed in patent application WO 02/38079 provides an improvement over the traditional pelvic floor repair surgery discussed above. WO 02/38079 discloses a specific mesh configuration for use in treating both cystocele and SUI that includes a basic mesh structure with two front and two rear supports that extend outwardly therefrom. The supports are thin strips of mesh that provide support to the prolapsed organ in place of sutures. In particular, WO 02/38079 describes using the instrument for treating SUI described in WO 96/06567 and 97/13465, which are incorporated herein by reference in their entirety, and the method described therein, to pass the mesh strips through the body. The described needles are used to pass the mesh strips from the vaginal incision up through the abdomen. One pass is made with a needle coupled to both the front and rear strips on one side, and then another pass is made for the front and rear strips on the other side. Thus, the front and rear strip on a given side are coupled to a needle together, and passed through the same channel out through the abdomen. As a result of both the front and back strips being attached at the same point, some additional stabilization is required to prevent forward movement of the basic structure, with suturing being the disclosed means by which to achieve this stabilization.\nAs with the traditional pelvic floor repair procedures, the procedure described in WO 02/38079 requires extensive, deep suturing skills in tight spaces, which are difficult to perform and time consuming. Sutures provide for only a “pin point” attachment of the mesh in place, and failures between the suture and the mesh are not uncommon. In addition, the single point of attachment of both mesh strips significantly limits the adjustability of urethral suspension versus bladder suspension. In particular, if tension free suspension of the urethra is achieved, further adjustment of the bladder suspension strips (rear strips) can adversely affect the urethral suspension, and vice versa. Finally, the vaginal approach of WO 02/38079 is disadvantageous in that suspension of the rear strip is limited to the same path as that of the front strip, as any other pathway would be potentially dangerous to vital organs and nerves via such a blind approach.\nAccordingly, it would be beneficial to provide an improved surgical system and method for pelvic floor repair."}
{"text": "Endoluminal prostheses are typically used to repair, replace, or otherwise correct a diseased or damaged blood vessel. A prosthesis may therefore be used to prevent or treat a wide variety of vascular ailments such as stenosis of the vessel, thrombosis, occlusion, or an aneurysm.\nOne type of well-known endoluminal prosthesis used in treatment and repair of diseases in various blood vessels is a stent. A stent is a generally longitudinal tubular device which is useful to open and support various lumens in the body. For example, stents may be used in the vascular system, urogenital tract and bile duct, as well as in a variety of other applications in the body. Endovascular stents have become widely used for the treatment of stenosis, strictures, and aneurysms in various blood vessels. These devices are implanted within the vessel to open and/or reinforce collapsing or partially occluded sections of the vessel.\nStents are generally open ended and are radially expandable between a generally unexpanded insertion diameter and an expanded implantation diameter which is greater than the unexpanded insertion diameter. Stents are often flexible in configuration, which allows them to be inserted through and conform to tortuous pathways in the blood vessel. The stent is generally inserted in a radially compressed state and expanded either through a self-expanding mechanism, or through the use of balloon catheters.\nA graft is another type of endoluminal prosthesis which is used to repair and replace various body vessels. Whereas a stent provides structural support to hold a damaged vessel open, a graft provides an artificial lumen through which blood may flow. Grafts are typically tubular devices which may be formed of a variety of materials, including textile and non-textile materials. Grafts also generally have an unexpanded insertion diameter and an expanded implantation diameter which is greater than the unexpanded diameter.\nIt is known to combine a stent and a graft to form a composite endoluminal prosthesis. Such a composite medical device provides additional support for blood flow through weakened sections of a blood vessel. In endovascular applications the use of a stent/graft combination is important because the combination not only effectively allows the passage of blood therethrough, but also ensures the implant will remain open.\nIt is also known to provide delivery systems for delivering grafts, stents and composite stent/graft prostheses intralumenally. These delivery systems generally include catheters with the prosthesis removably mounted to the distal end of the catheter. Quite often a catheter, introducer sheath, or other similar retaining means, is disposed over the prosthesis to removably support the prosthesis on the catheter. Once the prosthesis is situated in the target site in the lumen, the catheter is removed from the prosthesis.\nIn treating an abdominal aortic aneurism, traditional open surgery or minimally invasive endovascular procedures may be employed. Traditional open surgery requires a large incision in the abdominal wall, from just below the breast bone to the top of the pubic bone. The muscles are then divided and the intestines and internal organs of the abdomen are pulled aside the aorta in the back of the abdominal cavity, just in front of the spinal column.\nThe aorta is clamped and the aneurysm is cut open to reveal any plaque and clotted blood inside. Degenerative tissue is then removed. An aortic graft is then sewn to the healthy aortic tissue above and below the weakened area so that, when finished, it functions as a bridge for the blood flow.\nAfter the aortic graft has been sewn in place and all bleeding spots controlled, the aneurysm sack which has been opened along its length is sewn back up usually over the new graft. This prevents the new graft from rubbing against the intestines, which can damage the intestine wall. The entire procedure is fairly traumatic, is relatively high risk, and requires a long recovery period.\nThe endovascular procedure utilizes the endoluminal prostheses mentioned above to minimally invasively treat an abdominal aortic aneurysm. The endovascular procedure requires two small incisions in the groin. X-ray imaging typically guides the vascular graft, stent, or stent/graft through a blood vessel in the leg and into the aorta.\nWhile minimally-invasive techniques have been an improvement over prior open surgical techniques, there is still need for improvement. A drawback of systems to deliver endoluminal prostheses is the inability to adjust or retrieve the graft or stent/graft once it has been deployed. Generally, deployment of the graft (or stent/graft) marks a point of no return; if the graft is determined to be in an inappropriate position, or the graft size is inadequate, it is not always easy or possible to cancel the procedure or reposition the stent/graft after passing this point.\nAnother drawback of a known delivery system for delivering stent/grafts is that the large diameter introducer catheters are needed to deliver such systems. A typical previously known stent/graft system may include a central delivery shaft having a diameter of 1.5-1.75 mm, a deployment balloon having a thickness of 0.5-0.75 mm, an anchoring stent with a thickness of 0.3-0.6 mm, a synthetic graft with a thickness of 0.25-0.5 mm, and a delivery sheath having a thickness of 0.5-0.75 mm. A stacking of these thicknesses results in a combined thickness of 4-7 mm, which must be inserted through a vascular system generally having a diameter in the range of 5-7 mm.\nThis creates apparent problems in the delivery and trauma which may occur with such a prosthesis.\nU.S. Pat. No. 5,824,055 to Spiridigliozzi et al. is directed to providing a solution to many of these problems with the delivery systems. The delivery system employed by Spiridigliozzi et al. entails separate delivery of a stent and graft in order to provide a reduced profile and smaller diameter of the system. There are still, however, drawbacks to the system. In particular, once a graft is delivered first before the stent, because there is no anchoring system prior to the stent's introduction, the graft may migrate in the vascular system. This may result in imprecise location and increased trauma to the area. As such, this system requires the use of the stent to lock the graft in place.\nAccordingly, it is desirable to provide a system for treating vascular defects and disorders, such as abdominal aortic aneurysms, that minimizes the overall profile of the prosthesis and delivery system and results in a more secure graft/vessel interface."}
{"text": "1. Field of Invention\nThe present invention relates to integrated circuits and more particularly to modeling integrated circuits with multiple layers.\n2. Description of Related Art\nThe full-wave electromagnetic solver has become an essential part in the design flow for radio frequency integrated circuits (RFICs) serving as both a design and a verification tool. A three-dimensional broadband full-wave solver based on the integral equation (IE) method with the pre-corrected Fast Fourier Transform (FFT) acceleration has been developed to fulfill these functions [1].\nThe multilayer Green's function evaluation is essential to the performance of an IE solver. Various approaches have been developed to accelerate the efficient computation of multilayer Green's functions, such as the discrete complex image method (DCIM) [2], the fast Hankel transform (FHT) [3], and the window function based method [4]. The DCIM has been widely employed and improved since its debut. Two-level DCIM [5] was developed to split the original path into two. By sampling the paths at different rates, the method improved the efficiency and accuracy comparing the original one-level algorithm. Generalized surface wave pole extraction was developed to extend the DCIM to general multilayer media [6].\nFor on-chip applications, the semiconductor process must first be examined. The vertical cross-section of the metallization layers of a typical four-metal-layer semiconductor process is shown in FIG. 2. In order to fill the IE matrix entries efficiently, a 3D Green's function database is deployed [1]. The database is populated with the value of the Green's function between several source points at z′ and observation points at z for various horizontal distances ρ. Even with the aid of such a database, the rapid evaluation of the Green's function for each matrix entry is critical. Note that the distance between the adjacent layers is only a few microns or less, which causes the spectral kernel of Green's functions to decay slowly. Consequently, number of sampling points has to be increased to accurately capture the whole spectrum, which will slow down the DCIM dramatically because the generalized pencil of function (GPOF) method [9] in the DCIM has computation time proportional to O(Ns3), Ns is the number of sampling points. For that reason, accelerating convergence of the spectral kernel is necessary before applying the GPOF.\nThus, there is a need for improved methods for evaluating Green's functions for modeling multilayer integrated circuits."}
{"text": "Along with the advancement in computing technology, users' expectations of computational capabilities are similarly increasing. Users are constantly seeking resources that can provide the ability to achieve a computational result quickly and appropriately. Attending to users' requests is complicated by the fact that user projects vary in terms of required processing power, memory allocation, software capabilities, rights licensing, etc. Recently, systems have been organized to include a plurality of virtual machines. Tasks can then be assigned to virtual machines based on the task requirements, the machines' capabilities and the system load. However, given the dynamic nature of assignments and the many components in these systems, monitoring the systems' performance is difficult."}
{"text": "This invention relates to a new and distinct selection of Lagerstroemia indica, a member of the Lythraceae or Loosestrife family. Lagerstroemia indica cultivar Monink was discovered in a group of seedlings which originated from seed of Lagerstroemia indica `Little Chief`, seed was sown November 1983. During the summer of 1987, this new cultivar was selected at Monrovia Nurser Company, 18331 East Foothill Boulevard, Azusa, Calif. My new plant has been asexually reproduced by cuttings since 1987 at the above location in Azusa at Monrovia Nursery Company. The original seedlings displayed extreme variability, therefore the distinct phenotypic characteristics of my new selection that sets this plant apart from other Lagerstroemia indica plants would likely be lost through sexual reproduction. Therefore, sexual reproduction is prohibited and propagation is restricted to asexual reproduction by cuttings."}
{"text": "1. Field of the Invention\nThe present invention concerns an apparatus for protecting an electrical device from overtemperature, having a temperature-dependent switch with a housing on whose base a first connected contact is provided, a second connecting contact provided on the housing, and a holder which carries the housing and has an external terminal connected to the second connecting contact.\n2. Related Prior Art\nAn apparatus of the kind mentioned above is known from DE 38 17 080 C2.\nThe known apparatus comprises a holder made of plastic, in which a temperature-dependent switch with a two-part metal housing is arranged. The connecting contacts of the temperature-dependent switch are constituted on the one hand by the metal cover and on the other hand by the metal lower part, a bimetallic switching mechanism (known e.g. from DE-A-21 21 802) being arranged in the interior of the switch.\nThe known apparatus is slid onto connector lugs of an electric motor and itself has a connector lug which is accessible from outside.\nThe connector lug of the holder is connected to a thermal contact plate which, in the case of an apparatus placed on the electric motor, rests directly on the energizing winding. In the interior of the holder, the metal housing rests with its base on the contact plate, so that there is thermal contact with the electric motor, and electrical contact with the connector lug of the holder.\nPressing on the metal cover is a connecting spring which is joined to a connector lug of the electric motor but is no longer accessible from outside once the apparatus has been put in place. The second connector lug of the electric motor, and the connector lug of the holder, serve as external terminals of the electric motor equipped with the known apparatus.\nWith the arrangement described, the temperature-dependent switch is arranged in series with the electric motor between the two external terminals. If the temperature of the energizing winding increases above a defined value, the bimetallic switching mechanism then opens the electrical contact between the metal cover and metal lower part, so that an electrical voltage applied to the external terminals can no longer power the electric motor. This prevents the electric motor from heating up, during operation, above a defined temperature; i.e. it is protected from overtemperature.\nFor many applications the known apparatus is not suitable, since the holder used therein is on the one hand much too bulky to be mounted on small electrical household appliances, for example a coffee machine, an electric hotpot, or electrical hotplates. A further disadvantage of the known apparatus lies in the fact that it cannot track rapid temperature changes without delay, since heat is dissipated to the connector lugs via the metal cover and the metal lower part.\nRapid reaction to a temperature rise is, however, required specifically in the case of electric hotpots or hotplates, for example in order to prevent the hotpot, or a pot sitting on the hotplate, from boiling dry. Protection from boiling dry cannot be achieved by monitoring power consumption, since the amperage remains relatively constant regardless of whether or not water is still present; the temperature of the electrical device to be protected simply rises, rapidly and steeply, when the water has boiled off.\nApparatuses which protect household appliances from overtemperature are subject in general to particularly stringent demands in terms of electrical, mechanical, and thermal reliability, which are not met by the known apparatus."}
{"text": "A. Field\nThis invention relates to a thread clamp, in particular for a weaving machine, as well as to a weaving machine using such thread clamp.\nB. Related Art\nMore particularly, the invention relates to a thread clamp which at least consists of an actual clamp with at least two clamping parts between which a thread can be clamped, as well as of actuation means for moving said clamping parts in mutual respect between a closed position, whereby the clamping parts are pushed towards each other, and an open position whereby the clamping parts are situated located at a distance from each other.\nIt is known that such thread clamps can be used at different locations in a weaving machine. In particular, they are used for gripping weft threads, either for holding them at their place, or for braking and/or stopping them in their movement.\nIt is also known that such thread clamps can be actuated in different manners.\nSo, for example, electric actuating means can be used, whereby use is made of an electromagnet or a piezo-actuator for controlling the clamp, for which of course an electric wiring must be provided up to the thread clamp concerned, which often is viewed as disadvantageous, especially if such thread clamp is provided on a moving part, such as the sley of the weaving machine.\nIt is also known to apply pneumatic actuation means, whereby such thread clamp then is provided with a movable part which can be actuated by means of fluid (air) pressure.\nWith the known electric as well as the pneumatic embodiments, that is the electric drive and the pneumatic drive, respectively, the clamp is forced into one well-defined position, whereas the clamp, when switching off the actuation means, automatically is brought into the other position, by means of a return spring. The use of such return spring has various disadvantages. A first disadvantage consists in that it is rather difficult to mount a return spring which has a relatively large stroke and which does not vary too much in pre-tension, into a small built-in space, which return spring moreover, notwithstanding the relatively large displacement, does not buckle and also still has a long service life. A second disadvantage consists in that in certain applications, for example, when the thread clamp is mounted on a sley which is moving rapidly to and fro, it may occur that such return spring is excited automatically, for example, due to inertia forces, at moments when this is not desired. A third disadvantage consists in that mutual friction between the windings of the spring may occur, as well as between the spring and the surrounding housing, as a result of which the proper spring functioning can be disturbed. A fourth disadvantage consists in that, with high speeds, undesired effects can be created in such return spring, for example, as windings, so to speak, beat against each other and create undesired effects. A fifth disadvantage consists in that such return springs, due to their manifold movements, are subjected to intense wear and tear.\nIt is known to apply actuation means which are double-acting, for example, electromagnetic or pneumatic drives with which the thread clamp, by means of a controlled actuation, is forced into opened condition as well as into closed condition, without using passive readjustment means. This, however, has the disadvantage that a higher number of cables and compressed air conduits, respectively, have to be brought up to the thread clamps concerned, which in particular is a disadvantage with thread clamps which are mounted on a moving part.\nThe invention provides a thread clamp for a weaving machine which allows new constructional possibilities, control possibilities and combinations of controls, and is very efficient in its use and with which, depending on its form of embodiment, one or more of said disadvantages can be excluded or at least minimized.\nTo this aim, the invention in the first instance relates to a thread clamp for a weaving machine, which at least consists of an actual clamp with at least two clamping parts between which a thread can be clamped, as well as of actuation means for moving said clamping parts in mutual respect between at least two positions, a closed position, whereby the clamping parts are pushed towards each other, and an open position whereby the clamping parts are situated at a distance from each other, respectively, with as a characteristic that said actuation means comprise at least a pneumatically actuatable part for bringing the clamp into at least one of said two positions, which part also can be actuated by means of a negative pressure, more particularly vacuum.\nBy using a part which can be actuated by means of negative pressure, the advantage is obtained that new forms of constructions, as well as mounting possibilities are created. Another advantage consists in that a new actuation method is provided for, which, on one hand, enables an efficient control and, on the other hand, allows various new combinations with existing techniques.\nPreferably, the thread clamp according to the invention comprises a pneumatically actuatable part which enables the actual clamp, on one hand, to be brought into at least one position by means of at least a negative pressure, and on the other hand, enables the clamp to be forced into the other position by means of at least another manner of actuation. Thereby, a control is offered which has the advantage that two control possibilities are obtained, to wit a control by means of negative pressure, more particularly vacuum, and a control by means of another manner of actuation.\nPreferably, the pneumatically actuatable part is configurated such that the clamp is forced into opened condition by means of feeding the negative pressure, whereas the closed condition is brought about by means of the other manner of actuation. The pressure difference which can be created by means of a negative pressure, more particularly a vacuum, even if it should be possible to create an absolute vacuum, in different application thus is too small to provide for a good clamping force, and preferably the negative pressure, thus, as aforementioned, is used for performing the opening of the clamp, whereas for closing the clamp, an actuation then can be provided for in another manner.\nFor the aforementioned other manner of actuation, preferably use will be made of a medium which is at overpressure, for example, compressed air, such that, depending on the pressure of the available medium, thus larger forces can be generated.\nIn the most preferred form of embodiment, the actuatable part has a common actuation surface to which the negative pressure can be fed, and the overpressure can be exerted upon, respectively. By using a negative pressure as well as an overpressure, it is thus possible to create, with one of the same actuation surface, forces in two opposed directions and to effect movements of the actuatable part in two opposed directions.\nIn the aforementioned case, the thread clamp may have a single conduit for pneumatically actuating said actuatable part, which conduit is used in is common for actuating the clamp towards the closed condition as well as for actuating the clamp towards the open condition. An important advantage then consists in that only one operation conduit is necessary, in other words, only one pneumatic conduit, as a consequence of which the pneumatic circuit is considerably simplified, in particular at the location of the thread clamp itself. This advantage is particularly relevant when different thread clamps must be mounted next to each other or if such thread clamps must be mounted next to other parts, as then the small space available often poses a problem and imposes a limitation on the number pneumatic conduits locatable in the direct proximity of the thread clamp.\nThis is also particularly useful with thread clamps which are mounted on moving parts of a weaving machine, for example, on the sley, as then, thanks to the invention, the number of connecting conduits between the moving parts and the fixed control may remain limited.\nIn the most preferred form of embodiment, said actuation means are designed such that the movement of the actuatable part is performed exclusively pneumatical, with which it is meant that the movement of this part exclusively takes place under the influence of negative pressure and/or overpressure, realized by means of fluid, possibly, however, assisted by the gravity of certain parts, but surely not assisted by additional means for exerting a displacement force, such as return springs, electromagnetic drive means or the like.\nThanks to the fact that negative pressure is applied, as well as by combining this preferably with a control by means of overpressure, such additional means for exerting a displacement force in fact can completely be excluded, with the advantage that such thread clamp constructively can be realized in a particularly simple manner.\nIn a practical form of embodiment, said actuatable part consists of at least one small movable piston, the displacement of which controls the opening or closing of the clamp.\nPreferably, this small movable piston has a piston surface which opens to a chamber in which an overpressure and negative pressure, respectively, can be created, which permits the small piston to be actuated from one side for performing the displacement thereof in one direction, as well as in the opposite direction, by creating either an overpressure or a negative pressure, more particularly vacuum, in the aforementioned chamber.\nIn the most preferred form of embodiment, the small piston substantially will consist of a small plunger, one extremity of which forms one of the aforementioned two parts of the actual clamp, whereas the other extremity functions as a piston surface. Hereby, a particularly compact construction of the thread clamp is enabled.\nOf course, the invention also relates to weaving machines using at least one thread clamp as described in the aforegoing.\nImportant applications of such thread clamp on a weaving machine are as follows:                as a thread clamp at the inlet of the shed of the weaving machine;        as a thread clamp which cooperates with a main blower or auxiliary main blower;        as a thread clamp which cooperates with a main blower or auxiliary main blower which is attached on the sley of the weaving machine;        as a thread clamp at the end of the shed for holding the end of an inserted weft thread;        as a thread clamp on a moving part of the weaving machine, more particularly on the shed;        as a thread clamp which cooperates with a thread-preparing system.        "}
{"text": "Along with the improvement of laser beam performance (particularly, improvements in the technology for increasing the output of laser beams and the technology for shortening the wavelengths of laser beams), laser processing technology using lasers (e.g., using a laser as a light source for lithography in LSI manufacturing) has grown to be one of industries' core technologies. In addition to increasing the output of laser beams and improving the efficiency of laser beams, shortening of the wavelengths of laser beams is desired for implementing finer and more precise processing.\nIn recent years, nonlinear optical technology using high intensity attosecond light sources is developing (see, for example, Non-Patent Documents 1 and 2). Such optical technology is targeted for a soft X-ray wavelength range."}
{"text": "This invention relates to certain novel formulations of turf additives that act in such a manner as to permit proper amounts of moisture to contact root systems in order to reduce dry spots within highly managed turf areas and/or lawns. It is theorized that the accumulation of humic acid (and other natural byproduct deleterious compounds) at the topsoil surface in most cultivated grassy areas including sandy soils (such as, in particular, golf greens, pastures, lawns, and the like), as well as other non-grassy sandy areas in which such humic acid accumulation poses water repellency problems (such as beaches), results in the production of an effective organic waxy coating on the soil and/or sand components. Such a coating is hydrophobic in nature and thus dries out the soil itself. The inventive formulation thus permits removal of such accumulated humic acid (and other compounds) from the topsoil to the level necessary to provide effective moisture penetration for sustained grass growth therein (hereinafter referred to as xe2x80x9credistributesxe2x80x9d). Methods of providing such beneficial removal of bumic substances from target sandy soils are also contemplated within this invention, as well as specific test hydrophobic sand formulations.\nLocalized dry spots are a distinct problem within highly managed turf areas and/or lawns, in particular those with sandy soils, primarily for aesthetic reasons. Such dry spots are the result of the development of areas of varying degrees of water repellency within and at the surface of the target soil. Plant water usage is critical to sustained plant growth; however, the existence of such localized dry spots creates a problem with nonuniformity of water supply to treated grasses over time. Basically, in times of high stress and/or easy water evaporation (e.g., higher temperatures, low humidity), such water repellency areas will exhibit higher water loss than others. As a result, the plant life present within the target lawn or green will not receive uniform, and, at times, vastly different levels of, water supply. As time passes, the difference in the amount of water supplied to discrete areas of the target lawn or green may become more disparate. Thus, the possibility for localized dry spots to materialize within sandy soils is relatively high over a sustained length of time (e.g., from 6 to 18 months on average from genesis to being empirically noticed), and, again, most times the existence of such dry spots is unknown to the lawn or green caretaker until materialization (since the presence of such water repellency areas may exist anywhere within the topsoil, from the surface to as low as about 2 inches below, the area of greatest concentration of grass root systems).\nAlso, hydrophobicity of sand creates certain problems with regard to pooling water after raining (as one example) which in turn causes unsightly areas either within highly sandy yards, ballparks, or beaches, or to provide water penetration in dry sandy conditions in order to possibly sustain plant-life therein (such as arid desert-like areas). Reduction in such water repellency would thus be helpful in maintaining, at least, better aesthetics for such sandy areas, as well as the possibility for permitting or promoting the growth of sustained plant life in such dry, barren areas.\nWithout intending to be bound to any particular scientific theory, it is believed that such water repellency areas within sandy soils are the result of the presence of humic substances and their attachment to soil components, particularly in large accumulations at the topsoil surface. Humus is degraded plant and animal matter (by microbial organisms) and is basically the organic portion of soil that comprises the necessary nutrients to sustain plant growth and life therein. One byproduct of such humus (again produced through a naturally occurring process within the soil) is humic acid (simply the acidic form of humus, basically a mix of various different materials). Humic acid and other like substances, although necessary for the sustenance of plant life as it provides the aforementioned nutrients to root systems, unfortunately also appears to create problems within sandy soils, most particularly the creation of a waxy organic, water-repellent coating upon binding to and with soil components (for instance, and without limitation, sand). If such a coating is permitted to accumulate over a long period of time, such as the aforementioned 6 to 18 month period, and particularly at the topsoil surface, the coating becomes highly water repellent in nature and uniform plant water use is difficult to achieve, as mentioned previously. In theory, and, again, without intending to be bound to such theory, it is believed that such a coating is formed by the amphiphylic humic acid (or other like humic substance) adhering, by its hydrophilic portion, to the hydrophilic sites within the sandy soil, permitting the highly hydrophobic ends to extend (similar in nature to a micelle). Such a coating is thus hydrophobic in nature and, when present as a thorough coating over such surface portions, again, tends to either drive water away or facilitate water loss by preventing moisture from passing through to the subterranean roots of any plants therein. If the water remains at the surface, evaporation is also facilitated as it cannot easily penetrate within the soil. Such a problem exists, as noted above, not only within greens, but also within lawns and pastures (as merely some examples of such trouble areas). In order to provide a uniform appearance in lawns and greens, it has been a requirement either to water consistently in very large amounts (which is wasteful and possibly damaging to the plants themselves) or to water selected trouble areas by hand on a continuous basis (which is labor-intensive and possibly wasteful in terms of water consumption). Furthermore, it is generally too late to know of trouble water repellent areas within such lawns or greens until they become apparent empirically. For pastures, pools of water develop sporadically on occasion due to this problem; the standard method of remedying this problem is to dig up the earth and wait for the humic substances to be consumed as nutrients (over a relatively long period of time) by the root systems therein. Such a procedure thus leaves an aesthetically displeasing result and is not always reliable for reducing water repellency therein. Thus, it has been found that there exists a need to provide a simple method and/or formulation for reducing such a humic acid water repellency effect at the soil surface without creating detrimental effects to the surface plant life.\nIn the past, the best methods of reducing the amount and presence of localized dry spots have basically involved the introduction of certain surfactants to the soil for the transport of water through the surface coating, preferably in tandem with compounds that decrease the surface tension of the waxy coating to permit penetration of the active surfactant components themselves in U.S. Pat. No. 5,921,023 to Ogawa et al., U.S. Pat. No. 5,595,957 to Bowey et al., and U.S. Pat. No. 5,731,268 to Taguchi et al. Such a method has been problematic to a certain extent due to the cost associated with some silicon-based surfactants, biodegradability issues of most viable surfactants, as well as foaming problems when water is present, and/or the difficulty in removal of degraded coating components after surfactant treatment. Also, this specific surfactant-only treatment does not remove the waxy coating to an appreciable degree from the target topsoil surface.\nAnother manner of reducing such dry spot problems has been increasing watering itself. However, such a method is labor intensive and, in many areas where water is not plentiful, use for aesthetic purposes (e.g., lawns, greens, and the like), is preferably kept at a minimum as compared to other more important purposes (e.g., drinking water). Such an issue also contributes to the aforementioned development of water repellency areas over long periods of time because of the inability of the caretaker to continuously supply moisture to target lawns, greens, etc., to the levels needed to best ensure uniformity of watering is accomplished. Other possible attempts at alleviating such a problem exist, albeit as an aim at removing contaminants (e.g., oils, fuels, etc.) from the target soils for improving plant growth therein (U.S. Pat. No. 6,090,896 to Jahnke et al. and WO01/26832 to Lubrizol Corporation). Such a method does not provide the same degree humic substance removal, however, needed to provide the reliability of treatment for localized dry spot reduction. Also, another option has been the tearing up of the topsoil itself to physically breakdown the humic acid accumulation. Unfortunately, such a procedure is also labor intensive and invariably disturbs the plant life to too great an extent in order to properly manage uniform growth of target plants and grasses. Furthermore, once the plant life has begun to grow in such a scenario, humic acid begins to accumulate anyway, such that unless physically altering the topsoil is undertaken, the same problem of plant life loss (and thus the development of dry spots within the lawn or garden) would result without further effective moisture penetration. Thus, there remains a need for a simplified, chemical remedy to this waxy coating of humic acid (and other like substances) within lawns, greens, and the like. To date, other than those specific procedures noted above, such a remedy has not been provided with the aim at reducing such unwanted humic substance topsoil surface accumulation in order to provide sustained uniform water supply over the entire target green and/or lawn. Further previous attempts at remedying such localized dry spot water repellency problems include treatment of golf greens with sodium hydroxide. Such a treatment does in fact remove the hydrophobic coating; however, it is also highly phytotoxic to grass and thus is an unacceptable method from a commercial perspective.\nIt is thus an object of the invention to provide an improved method of lowering topsoil surface tension by chemically reacting (or possibly complexing) with the hydrophobic portion of such humic acid coatings, thereby permitting moisture to penetrate such a coating and not only pass to the root systems of target plants therein but also possibly transporting the surface-accumulated humic substances into the soil for root consumption of the nutrients provided by such substances. A simple, safe formulation permitting such a method is also an object of this invention. Additionally, a test hydrophobic sand formulation is also an object of this invention in order to permit reliable laboratory analysis of the effectiveness of certain formulations to provide the aforementioned desired reaction with hydrophobic portions of humic acid coatings.\nAccordingly, this invention concerns a soil additive formulation and/or method of treating sandy areas, soils, or areas including both sand and soil (such as lawns, greens, pastures, beaches, dry desert-like areas, and the like), wherein said soil additive formulation is non-phytotoxic and exhibits a humic acid removal capacity under the hydrophobic sand humic acid removal test of at least 150 ppm. Such a formulation comprises, preferably, at least one humic acid redistribution compound selected from the group consisting of at least one C2-60 (or possibly higher) alkyl, alkenyl, and/or alkaryl succinic anhydride or acid (the anhydride will convert to its acid form upon dissolution in water) mono-cation or di-cation salt resulting from the reaction with either a metal or amine functional neutralization salt, at least one C4-60 (or possibly higher) branched or unbranched diacid mono- or di-cation salt resulting from the reaction with either a metal or amine functional neutralization salt, at least one polycarboxylic acid derivative salt, and any mixtures thereof, and from 0-99% by weight of at least one compound that actively lowers the surface tension of humic acid waxy coatings from hydrophobic sand particles. Such a formulation may also comprise a copolymer exhibiting both hydrophilic and hydrophobic portions for reaction with the hydrophobic portions of such hydrophobic sand particles in order to further provide hydrophilic extensions therefrom to facilitate the reaction between the bound humic acid and the aforementioned humic acid removal compound. Furthermore, an inventive synthetic hydrophobic sand formulation exhibiting a penetration period of at least 9 seconds for a drop of 2 molar ethanol under a minimal ethanol drop test is also encompassed within this invention. Lastly, a method for reducing localized dry spot formation within lawns or greens comprising the application of a soil additive formulation to a target lawn or green, wherein said soil additive formulation exhibits a humic acid removal capacity under certain conditions of at least 150 ppm from a sample of the aforementioned synthetic hydrophobic sand formulation is additionally encompassed by this invention. Again, to date, nothing within the pertinent prior art teaches or fairly suggests such specific inventions.\nSuch a composition and method of treating sandy areas may thus be utilized for the redistribution of humic substances in sandy areas alone. In such a manner, the sandy area (if it a beach, for example) may be modified to permit water penetration therein, to prevent unsightly water pools, for example, after raining, or to dry desert-like areas in order to permit water penetration to sustain root systems of plant-life which would not grow otherwise.\nThe term xe2x80x9csyntheticxe2x80x9d as it applies to the inventive hydrophobic sand above, is intended to indicate that such sand is manipulated via a synthetic route to provide such a high level of hydrophobicity through chemical treatment of sand itself. It is not intended to mean a sand formulation that has been produced in total by hand. Of course, sand, being silicon dioxide, is available naturally, but it generally exhibits very low, if any hydrophobicity levels. Thus, the term xe2x80x9csyntheticxe2x80x9d is, again, intended to show that the hydrophobic properties thereof are chemically provided by hand and not in nature.\nThe inventive formulation may either be applied in liquid form, pellet form, or granular form to the selected treated area.\nThe inventive formulation, in terms of composition, thus requires at least one humic substance (in particular, humic acid) redistribution compound. In effect, such a compound or mix of compounds reacts with the hydrophobic portions of the humic acid (the portion of the humic acid not attached to the hydrophilic sand and/or soil particles) present at the target topsoil surface. Upon binding thereto, the redistribution compound provides sites for strong water adhesion. Such adhered water droplets will be pulled into the sand and/or soil by further adhesion by other particles or through cohesion with other water droplets). The binding energy of the water droplets to the redistribution compound is greater than that of the humic acid to the sand and/or soil, thereby permitting release of the humic acid for transport into the target soil for consumption as nutrients by the root systems of the plants therein. Thus, the humic acid redistribution is effectuated sufficiently to allow for greater amounts of moisture to penetrate the topsoil as well as effectively permitting transport of humic substance nutrients to the target roots, all for increased plant growth and reduction of water repellent areas.\nSuch a humic acid redistribution compound may be of any type that provides the above-discussed humic acid redistribution and transport via water movement. Preferably, and without limitation, such a redistribution compound may be chosen from three different classes of salts or weak polymeric acids, namely succinic anhydride (or acid) mono- and di-cation salts resulting from reactions with a metal or amine functional neutralization salt, such as hydroxy, alkoxy or C1-C60 alkyl, alkenyl, and/or alkylaryl succinic acid mono- and di-cation salts resulting from reaction with a metal or amine functional neutralization salt; C4-C60 branched or unbranched diacid salts, and polycarboxylic acid derivative salts.\nNon-limiting, preferred compounds for this purpose include alkenyl succinic anhydride di-cation salts, in particular such compounds as octenyl succinic anhydride and tetrapropenyl succininc anhydride di potassium salts, available from Milliken and Company under the tradenames SYNFAC(copyright) 8515 and SYNFAC(copyright) 8510, respectively. Examples of non-limiting preferred diacid salts within this invention include potassium sebacate, potassium adipate, and other dicarboxylic acid mono-metal, di-metal, or amine functional mono- or di-cation salts including such basic compounds as oxalic, malonic, succinic, glutaric, tartaric acid, and malic acids, having metal or amine counter ions of any corresponding valence metal or neutralizing amine, such as, without limitation, Na, Li, K, Mg, Ca, monoethanolamine, diethanolamine, triethanolamine, ammonia, and/or monoalkyl, dialkyl, and/or trialkyl amines. Such a compound (or compounds) is one that can be easily prepared by those skilled in the art by neutralizing a selected dicarboxylic acid with a selected base to form the desired salt or amine. Examples of preferred, non-limiting, polycarboxylic acid salts include solutions of maleic copolymer salts, water soluble polymeric polyelectrolytes, and other sodium salt-based anionic polymers, available from Milliken and Company under the tradename INVIGORATE(copyright) (a formulation previously utilized for the purpose of agglomerating fine soil particles into larger particles thus creating greater pore space therebetween as well to agglomerate clay or organic fines). Each of these classes of humic acid redistribution compounds provides the requisite humic acid removal and transport discussed previously. As noted below, such compounds effectively remove such unwanted topsoil substances from synthetic hydrophobic sand and thus are effective as compounds utilized by themselves for such purpose when applied to target lawns and/or greens (as one preferred embodiment). The humic acid redistribution compound may thus be comprised of all of the soil (and/or turf) additive formulation, but such a formulation preferably comprises from 1-99% by weight of such a redistribution compound; more preferably from about 10-90% by weight; more preferably from about 15-75% by weight. For the succinic anhydride-type salts, the amount may be even more preferably from about 50-75% by weight of the formulation, still more preferably from about 60-72%, and most preferably between about 65-72%. For the maleic copolymer alternative compounds, the amount may be even more preferably from about 10-25% by weight of the entire formulation, still more preferably from about 12-20%, and most preferably between about 17-19%. In such an instance, however, a large amount of the maleic copolymer-containing formulation will be water, from about 40-60% by weight, more preferably from about 40-50%, in order to permit uniform dispersion for effective application to the target lawn and/or green.\nHowever, in order to best ensure initial penetration of such removal compounds within the target topsoil areas (which may or may not be thoroughly coated with humic substance waxy coatings), it is preferable to include at least one compound within the formulation for the lowering of the surface tension at the topsoil surface which is also compatible with the aforementioned required humic acid removal compound. Such a compound can be an alkoxylated (preferably ethoxylated) alcohol (surfactant), such as a branched or unbranched C6-C60 alcohol alkoxylate (preferably, again, ethoxylate) (for utilization with the aforementioned succinic anhydride salts and long-chain acid salts), or alkoxylated (preferably ethoxylated) C8-C40 fatty acid (for utilization in combination with the aforementioned maleic copolymer-containing formulations). Such compounds may be branched or unbranched in configuration. Examples of preferred types of alcohol alkoxylates for this purpose include C6-60 alkyl,alkenyl or alkylaryl EO/PO surfactants, linear or branched, and secondary or primary hydroxyl in type, including mixtures of surfactants comprising from 99 to 1% by weight of at least one surfactant selected from polyalkylene oxide compounds with the general formula:\nR3xe2x80x94Oxe2x80x94(C2H4O)c(C3H6O)dR3\nwherein c is 0 to 500; d is 0 to 500, and R3 is H, or an alkyl group with 1 to 4 carbon atoms; wherein the polyalkylene oxide has a molecular weight in the range of 300 to 51,000; and a second optional different surfactant comprising a compound of the general formula\nR4xe2x80x94O(CH2CH2O)x(CHR5CH2O)yR6\nwherein x is from 1 to 50; y is 0-50: R4 is a branched or linear alkyl, alkenyl, aryl or an aryl group optionally having an alkyl group substituent, the alkyl group having up to 60 carbon atoms; R5 is selected from H and alkyl groups having from 1 to 2 carbon atoms; and R6 is selected from H and alkyl groups having from 1 to 30 carbon atoms. Suitable secondary surfactants also include carboxylic and dicarboxylic esters of the general formula:\nxe2x80x83R4COa(CH2CH2O)x(CHR5CH2O)yCObR6\nwherein x is from 1 to 50; y is 1-50, a is from 1 to 2, b is from 1 to 2: R4 is an alkyl or alkenyl group having up to 60 carbons or an aryl group optionally having an alkyl group substituent, the alkyl group having up to 60 carbon atoms; R5 is selected from H and alkyl groups having from 1 to 2 carbon atoms; and R6 is selected from H and alkyl groups having from 1 to 30 carbon atoms.\nThe surface tension of such a surface-active compound (or compounds) should in effect be below the general level of such a humic substance waxy coating, thus less than about 30 dynes/m2. As non-limiting examples for this purpose, tridecyl alcohol (8 EO), and coconut fatty acid (9 EO), are preferred. The amount of such a component within the inventive formulation is, as above, different for each type, depending on the type of humic acid removal compound present therein. Thus, for the succinic anhydride and long chain salt types, the amount of such an additive should range from about 1-20%, more preferably from about 5-15%, and most preferably from about 7-10%, all by weight of the entire formulation. For the remaining type of humic acid removal compounds, the amount of such an alkoxylated fatty acid ranges from about 1-25%, more preferably from about 10-20%, and most preferably from about 16-20%, all by weight of the entire formulation. Such an alkoxylated fatty acid is essentially required upon the presence of such maleic copolymer type components, if such a type of removal compound is actually present therein.\nAnother optional compound for introduction within the inventive formulation is a polyoxyalkylenated copolymer additive comprising at least two different alkylene oxide monomers, such as, without limitation, but preferably, ethylene oxide and propylene oxide, may be added, particularly with the highly preferred aforementioned succinic acid salts and/or long chain salts, for facilitation of the binding of the redistribution compounds to the target humic acid deposits through such copolymer groups. In essence, without intending on being bound to any scientific theory, it is believed that the more hydrophobic portions (propylene oxide, or PO, monomers, for example) bind to the hydrophobic ends of the humic acid, and the more hydrophilic portions (ethylene oxide, or EO, for example) bind readily to the salt removal compounds. In such a manner, the binding energy of the humic acid to the soil and/or sand is overcome by the pull of water droplets on the entire complex of salt-copolymer-humic acid such that transport to within the soil is accomplished and the humic acid can thus be more readily removed from the topsoil surface. Again, such a copolymer component is not necessary for proper functioning of the inventive formulation in every instance, although its presence may be desired in an effort to reduce the amount of humic acid removal compounds (which may be expensive or difficult to find in large quantities) within the soil (and/or turf) additive formulation and still provide an effective manner of reducing localized dry spots within the target lawn and/or green. Such a copolymer may thus be of any length and molecular weight with a preferred molecular weight of between 1000 and 5000, more preferably from about 2000 to about 3500, and most preferably from about 2750 to about 3250. Such a copolymer is available from BASF under the family of tradenames of PLURONIC(copyright). If present, such copolymer should be present in an amount of from about 1 to about 85% by weight of the entire formulation, more preferably from about 20 to about 80%, and most preferably from about 55 to about 75%.\nSuch an inventive formulation is one example of a soil additive that provides the desired topsoil humic acid removal that is necessary to effectuate the reduction in dry spot formation within vegetative areas. As discussed above, the aim of this invention is to lower the surface tension of humic acid accumulations on topsoil by reacting with the hydrophobic extensions of such an acid coating. This result is analogous to certain laboratory analyses involving the removal of humic acid from a synthetic hydrophobic humic-acid coated sand formulation. In order to best develop and test the effectiveness of such a novel soil additive formulation, it was necessary to develop the aforementioned novel synthetic hydrophobic sand since no such sand is available commercially. As noted below in greater detail, the basic premise in producing such novel sand is to attach large amounts of humic acid thereto in order to provide a highly hydrophobic sand. Such a sand formulation can be utilized for more than just this analysis of humic acid removal, as low water uptake can help in terms of less water-laden sandbox components, better flowing hourglass sand, and countless other end-uses. For this invention, the humic acid-treated sand exhibits a very high hydrophobicity in terms of adding drops of various concentrations of ethanol to the target sand surface (when flattened as a level surface) and determining the molar concentration of methanol required to permit penetration of the ethanol into the sand itself within a ten second interval. Again, as measured below, such an inventive sand must exhibit a penetration time of at least 9 seconds for a drop (e.g., about 40 {circle around (3)}L) of 2M ethanol (based on pure 200 proof ethanol). Upon producing this test sand, then, a proper laboratory analysis of the humic acid removal capability of the inventive soil additive formulation can be undertaken. Such a method of analyzing humic acid removal from such synthetic hydrophobic sand basically entails, for our purposes, the analyses of water-removed humic acid from hydrophobic (humic acid-coated) sand particles, as discussed in greater detail below.\nThe inventive formulations may include any other standard components for lawn, garden, or other vegetation treatment, including, further wetting agents, colorants (for aesthetic purposes or for application identification), perfumes, water, electrolytes, fertilizer, pesticides, and the like.\nInitially, the production of a suitable hydrophobic sand formulation was necessary to properly analyze the ability of the particular soil additive formulations for sufficient removal of humic acid coatings applied thereto. As noted above, such an hydrophobic sand formulation is novel as to its synthetic nature and its resultant reproducible hydrophobicity as measured by an ethanol drop test. Thus, such a novel hydrophobic sand formulation was produced and tested as follows:\nFirst, humic acid sodium salt was converted to acid form by treating with a molar amount of 1N hydrochloric acid followed by coating target sand with such humic acid. Thus, humic acid sodium salt (purchased from Aldrich Chemical) was mixed with 1N hydrochloric acid (36% acid diluted with Deionized water). Using 1 liter of such 1N HCl, 100 grams of Humic Acid sodium salt was added and allowed to stir 16 hours. After stirring, the solution was filtered through a Whatman #2 filter paper and washed with excess deionized water three times to remove excess HCl. The residual humic acid was then dried in an 80xc2x0 C. oven for 12 hours (or until dry). The resultant humic acid material was then ground with a mortar and pestle to a fine powder consistency.\nSand, a GA-45 washed top dressing sand for golf greens, available from Golf Agronomics, was then treated with such humic acid. Using 20 grams of the treated humic acid from above, it was combined with 250 grams reagent methanol and 250 grams deionized water. The solution was then agitated for two hours to insure full dissolution of the powder. Once dissolved, 500 grams of the sand was added thereto and the mixture was allowed to tumble for at least 4 hours to insure full coating by the humic acid. Subsequently, the sand was then filtered through Whatman #2 filter paper to remove excess solvents and then placed in an 80xc2x0 C. oven to dry for at least 16 hours. The treated sand was then removed from the oven and washed 3 times with deionized water through a Whatman #2 filter. The filtered sand was then placed in the same oven for another 16 hours or until dry.\nThe sand was then characterized by a minimum ethanol drop method. First, ethanol standards were made for use in the MED(Minimum Ethanol Drop) test through the production of 1M, 2M, 3M, and 4M solutions of ethanol using absolute 200 proof ethanol. A 15 mm petri dish with one eighth of an inch of the test sand was used for the MED test. Ten drops of distilled water were placed on top of the test sand and a stopwatch was used to record the penetration time. After five minutes, the drops were removed. Ten drops of the one molar ethanol were then placed on the sand and timed (an average of two minutes and 10 seconds). Ten drops of the two molar ethanol were then tested in the same manner (an average of nine seconds). This test required that the drops that last an average of ten seconds be given the numerical value of the molar solution tested. Thus, the produced novel synthetic hydrophobic sand formulation exhibited a MED for 2M ethanol of at least 9 seconds.\nSoil additive formulations were then produced for measurement in terms of humic acid removal from this test novel synthetic hydrophobic sand (all percentages listed below are byweight of the entire formulation):"}
{"text": "In the field of orthodontics, it is known to adhere orthodontic appliances, such as brackets, buccal tubes and the like, directly to teeth. Typically, this is accomplished by chemically bonding an appliance to a tooth surface using an adhesive. If desired, the bonding surface area of the appliance may be roughened in order to increase the surface area of the appliance in contact with the adhesive, thereby enhancing the chemical bond. Direct bonding also may be accomplished by using a metal appliance having undercuts for forming a mechanical bond with the adhesive.\nU.S. Pat. No. 5,267,854 to Schmitt teaches a metal injection molded orthodontic bracket including a plurality of raised posts extending bucco-lingually from a tooth abutting surface. Each post includes a root section having a base integrally formed with the tooth abutting surface and an apex section bucco-lingually extending from the root section. The apex section terminates in a sharp, continuous parameter edge that was originally smaller in all directions than the root section. However, in accordance with the Schmitt patent, further cold working of the parameter edges occurs such that each edge is worked at ambient temperatures into a mushroom-shaped button having a worked edge larger in all directions than its associated root section and an eave capable of mechanically bonding with an adhesive.\nU.S. Pat. No. 4,661,059 to Kanno features a metal orthodontic bracket which has a base surface provided with a plurality of orthogonal grooves formed by a cutting machine having a plurality of rotatable thin circular cutter blades. The grooves so produced have small fins or flashes at their edges resulting from cutting operations of cutter blades scraping metal matrix at high speed. These small fins are pressed down into the inside of the grooves to form undercuts for the adhesive to provide mechanical bonding of the base to the tooth surface.\nAlthough mechanical bonding provides some advantages over traditional bonding methods, the brackets taught by Schmitt and Kanno have several limitations. For example, both brackets are made of metal and therefore, lack the aesthetic qualities found in plastic or ceramic orthodontic appliances. Furthermore, the mechanical bonding surfaces of the Schmitt and Kanno brackets are formed by cold working the metal at ambient temperature. This cold working process also distorts the metal microstructure of the raised posts used for mechanical bonding.\nIn addition, the methods used to form the mechanical bonding surfaces are relatively complex and expensive. The Schmitt patent requires the use of a hydraulically activated metal coining die array to cold work the posts, with one die for each post, and with each die having a partially spherical concave mushroom forming cavity in the tip. Such die arrays are difficult to make, expensive to manufacture, and difficult to use. Kanno teaches the use of a cutting machine to form the posts in the bracket base, which leads to a significant amount of wasted metal. Also, the raised posts of the Schmitt bracket must be originally molded with an apex section narrower than the root section in order to remove the metallic bracket from the mold without significant risk of peg breakage. This tapering results in less metal being available at the apex to form the mushroom-shaped button and corresponding eave required for mechanical bonding. Furthermore, the extremely fine fins or flashes of the Kanno bracket formed by the cutting process must be pressed downward in order to form the mechanical bonding surfaces and are subject to stress and fracture in this process.\nIn the related patents and applications referred to above, the bracket bases are formed of plastic having arrays of pegs projecting from their bases, with the outer ends or tips of the pegs having mushroom shaped heads that present overhangs which facilitate mechanical bonding of the bases to the surfaces of teeth with adhesive. The heads are formed by heat softening the peg tips and, while softened, applying a force to the tips in a direction toward the base, deforming the tips into mushroom-shaped heads.\nTherefore, it would be beneficial to have an improved method of deforming peg tips of projections on orthodontic appliance bases into mushroom-shaped heads without the need to use heated platens to soften the projections to facilitate their deformation or to use expensive cold working coining dies in order to make an orthodontic appliance capable of being mechanically bonded to a tooth, in which the appliance is made of an aesthetically pleasing material."}
{"text": "Interferometric modulators are spatial light modulators typically used in display applications. These modulators use interference effects to control the intensity and the color of the light seen by a viewer. One example of such a modulator is the iMoD™. The iMoD employs a cavity having at least one movable or deflectable wall. As the wall, typically comprised at least partly of metal, moves towards a front surface of the cavity, interference occurs that affects the color of light viewed at the front surface. The front surface is typically the surface where the image seen by the viewer appears, as the iMoD is a direct-view device.\nThe front surface is transparent, as it is the portion of the device through which the viewer views the image. Manufacture of an interferometric modulator is somewhat similar to other types of devices manufactured on glass, such as liquid crystal devices (LCD). However, due to the unique nature of the device, packaging of the device has some other considerations.\nInterferometric modulators are microelectromechanical systems (MEMS) devices. While they use a glass substrate like LCDs, they also have moving parts and must have room to move while remaining protected. Similarly, while they are MEMS devices, they are manufactured on a substrate that has different processing considerations than most MEMS devices.\nThere are some degrees of freedom available to interferometric modulators not available to either LCDs or many MEMS devices. Unlike LCDs, interferometric modulators do not have material injected into the package, such as the liquid crystalline material is into LCD display panels. Unlike many MEMS devices, interferometric modulators do not have parts that rub against one another, alleviating concerns about internal particle generation."}
{"text": "1. Field\nThis relates to a dryer, and more particularly, to a dryer having enhanced dehumidifying power.\n2. Background\nIn a laundry treating apparatus having a drying function such as a washer or dryer, once washing and dehydration are completed, hot air may be supplied into the drum to evaporate moisture from the laundry, thereby drying the laundry. Such a dryer may include a drum rotatably provided within a cabinet, a drive motor to drive the drum, a blower fan to blow air into the drum, and a heating device to heat air conveyed into the drum. The heating device may use, for example, high-temperature electric resistance heat generated using electric resistance, or combustion heat generated by combusting gas."}
{"text": "This invention relates in general to apparatus for fusing toner images to toner image carrying supports through the application of heat and pressure. More particularly, this invention relates to fuser apparatus including a sensor for sensing the temperature of a heated fuser member and for producing a signal representative of the temperature of the member which is used to maintain the temperature of the fuser apparatus within predetermined limits and to indicate a jam condition when a sudden large drop in temperature is sensed.\nElectrographic copiers produce copies having toner images which are permanently fixed to supports such as copy sheets by well-known fusing techniques. In one such technique, a copy sheet carrying a toner image is contacted by a heated fuser member such as a roller or belt to permanently fuse the toner image to the copy sheet through the application of heat and pressure. The fuser member may be heated by means of an internal or an external heat source. In either case, the temperature of the fuser member is maintained within predetermined limits by sensing the temperature of the heated member and providing a feedback to control the heat source.\nThus, sensors for sensing the temperature of a heated fuser roller may be provided in contact with the fuser roller as disclosed in U.S. Pat. No. 3,291,466 for \"Xerographic Fixing Device\", issued Dec. 13, 1966, patentees G. A. Aser et al (see, for example FIG. 6, element 825; column 6, lines 56-68); U.S. Pat. No. 3,324,791, for \"Xerographic Roller Fuser Drive Apparatus\", issued June 13, 1967, patentees J. R. Cassano et al (see, for example, FIG. 2, element THS-2; column 5, lines 66-72); and U.S. Pat. No. 3,327,096 for \"Temperature Control Circuit\", issued June 20, 1967, patentee T. Bernous. Contact temperature sensors, although useful in controlling the temperature of the fuser roller, may cause damage to the fuser roller surface especially where the surface is of an elastomeric material. Copy sheets which wrap around the fuser roller may engage the sensor resulting in damage to the sensor and jamming of sheets in the fuser.\nTemperature sensors for heated fuser rollers may also be spaced a distance from the fusing surface thus obviating the difficulties of sensors which contact the fuser surface. Such sensors are disclosed in U.S. Pat. No. 4,297,562 for \"Heat Roller Type Fixing Apparatus\", issued Oct. 27, 1981, patentees N. Kamogawa et al (see, for example, FIGS. 1 and 2, column 1, lines 62-66); Research Disclosure No. 1370, entitled \"Electrophotographic Fuser Apparatus\", dated September 1975 (see, for example, FIG. 3 element 79 and description thereof); U.S. Pat. No. 3,369,106 for \"Process Heating Control System\", issued Feb. 13, 1968, patentee J. H. Troll (see, for example, FIG. 2, element 6, column 2, line 54 et seq.); U.S. Pat. No. 3,742,191, for \"Infrared Temperature Sensor and Control for Use With Heated Moving Bodies\", issued June 26, 1973, patentees R. R. Poole et al (see, for example, FIG. 6, element 51, column 4, line 52 et seq.). The latter two patents disclose temperature sensors for measuring infrared radiation from a heated roller. Instead of sensing temperature at the fuser roller surface, the temperature of the fuser core may be sensed to provide temperature control temperature of the fusing roller. Such core sensing controls are disclosed in U.S. Pat. No. 4,046,990 for \"Temperature Sensing and Control of a Fusing Roll\", issued Sept. 6, 1977, patentee W. E. White (see, for example, FIGS. 1 and 2, element 5, column 4, line 22 et seq.) and Research Disclosure No. 19642 entitled \"Temperature Control of Fuser Roller\" published August 1980. Where the fuser includes a heater element positioned over a copy sheet path, a sensor may be positioned on the other side of the path to detect temperature variation due to the passage of the copy sheet between the radiation source and the sensor such as disclosed in Japanese Patent Publication No. 76/22675 published June 11, 1976 entitled \"Fixing Device for Electronic Copying Machine\" (see for example FIG. 1).\nDuring the fusing process there is a tendency for a fused copy sheet to adhere to the fuser roller due to the stickiness of the toner material. Efforts to minimize sticking of copy sheets to fuser rollers include making the surface of the fuser roller of a material which has good release properties, such as silicone elastomer or polytetrafluoroethylene. Silicon release oil may also be applied to the fuser roller surface to enhance the release properties thereof. Since, however, there is still a possibility of a copy sheet adhering to the fuser roller, in commercial copiers a detector is provided in the copy sheet path near the exit nip of the rollers to detect whether or not a copy sheet has advanced past the nip. Such a detector for example is disclosed in U.S. Pat. No. 4,391,509 entitled \"Roller Fuser Apparatus In Which Copy Sheet Jams are Minimized,\" issued July 5, 1983, by W. A. Cavagnaro (see for example FIG. 2 element 70, column 5 lines 47-50 and column 6 line 15 et seq).\nThus, the control of the temperature of a heated fuser member and the detection of paper jams in fuser apparatus are desirable in present day copiers to minimize down time of the equipment due to malfunctions. With higher speed machines, there is also a trend to more complex operational and diagnostic functions. This, it would be desirable to provide copier controls which would simplify the control function while providing efficient and rapid system response to a machine malfunction such as a copy sheet sticking to a fuser roller. It would also be desirable to provide controls for a fuser roller which do not interfere with the fuser components and which do not contribute to fuser roller damage or degradation."}
{"text": "Wireless communication networks support communications between mobile communication devices. Such mobile communication devices are often referred to as subscriber units (SUs) and, in some cases, portable subscriber units (PSUs) or mobile subscriber units (MSUs). In some wireless communication networks, for example, land mobile radio (LMR) and trunked radio networks, the subscriber units can communicate with each other in talk groups. A talk group is a group of subscriber units, created by an administrator, in which each subscriber unit in the talk group may participate in a group call. Group calls may be initiated and managed, for example, at a call controller, a call center, a control head, or a server.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of embodiments of the present invention.\nThe apparatus and method components have been represented where appropriate by conventional symbols in the drawings, showing only those specific details that are pertinent to understanding the embodiments of the present invention so as not to obscure the disclosure with details that will be readily apparent to those of ordinary skill in the art having the benefit of the description herein."}
{"text": "Computing devices are conventionally configured in a variety of specific ways to support different user interactions with the computing device tailored for particular settings. The computing device, for instance, may be configured to support mobile usage and therefore is optimized for that usage, such as by including a battery, portable size, and user interface configured for interaction by one or more hands of a user. Likewise, the computing device may be configured for non-mobile applications, such as a traditional desktop PC that has a relatively large size and includes devices that are configured to aid user interaction in a desktop setting, such as keyboards and cursor-control devices.\nComputing devices are now available that support use in a variety of settings and consequently may support interaction in a variety of usage scenarios. Examples of such computing devices may include tablets and other configurations that include significant computational resources that enable these devices suitable to perform a variety of different interactions. However, conventional techniques typically limited a user to user interactions that were tailored for a single setting and corresponding usage scenario, and thus could limit interaction with the device when used in other settings."}
{"text": "Electronic card connectors of electronic devices are installed directly on a circuit board in the early, so a housing of the electronic device must be removed firstly to place an electronic card in the electronic card connector. While electronic card connectors of some electronic devices, such as mobile phones, are installed under batteries, in this case, the housing and the battery of the mobile phone all need to be removed sequentially, and then the electronic card can be placed in the electronic card connector, said removing and placing will bring a lot of trouble to consumers.\nIn order to comply with the development trend of simplifying electronic device, a drawer card connector appears recently, it usually includes an insulation body, multiple conductive terminals positioned at the insulation body, and a covering housing masking the insulation body. The insulation body and the covering housing coordinate with each other so as to form an accommodating space therebetween, and a sliding holder is accommodated in the accommodating space. An exiting mechanism is fixed on the insulation body for sliding the sliding holder forward and backward in the accommodating space. The exiting mechanism further includes a pop-up element, a pivoting element and a pushing rod, wherein, the pop-up element and the pivoting element are located at back of the insulation body, and at least part of the pop-up element exposes in the accommodating space and presses against the sliding holder elastically. The pivoting element has an actuating portion and a pushing portion, the pushing rod extends forward and backward at a side of the insulation body, a back end of the pushing rod presses against the pushing portion, a leading end of the pushing rod is extended to the front of the accommodating space.\nThe electronic device has an accommodating cavity defined at a side thereof, and the accommodating cavity has an opening communicating with outside, the drawer card connector is installed in the accommodating cavity, the electronic card can be placed in the drawer card connector without removing the housing of the electronic device. Just need pushing the leading end of the pushing rod with a tool that firstly enters into the opening, and the leading end of the pushing rod drives the back end of the pushing rod to press against the pushing portion backward, and then the pushing portion drives the actuating portion to rotate and thereby press against the sliding holder forward. Meanwhile, the pop-up element applies a forward thrust to the sliding holder, thereby the sliding holder exits outward, and then the electronic card can be put in the sliding holder, finally the sliding holder is pressed backward into the accommodating space, and the sliding holder can be positioned therein.\nThough the exiting function of the drawer card connector is easy, which can overcome the removing and placing problem of the conventional electronic card connector, but it doesn't comply with the development trend of miniaturization and low cost for electronic devices nowadays, that is because the exiting mechanism includes three separate elements, the pop-up element, the pivoting element and the pushing rod, which are disposed at back of and a side of the insulation body respectively, and occupy much space of the insulation body, thereby prevent miniaturization progress of the drawer card connector. In addition, the processes of manufacturing and assembly are complex, which increase the cost of manufacturing, thereby fail to achieve the goal of low cost.\nHence, it is desired to provide a new card connector to overcome the above-mentioned drawbacks."}
{"text": "During the processing of semiconductor wafers, wafer misalignment and wafer movement after placement on a wafer holder (e.g., susceptor) can cause the wafer to be outside the holder's wafer pocket, which is designed to hold the wafer in the correct position as the wafer is rotated. If a wafer is not fully seated in the pocket, the resulting temperature gradients across the wafer can make it unusable. Thus, it is of great importance to detect whether wafers are being processed while they are misplaced or out-of-pocket, in order to take corrective action before the wafer and/or equipment are damaged. In addition, due to the extremely harsh environment inside the process chamber, simplicity and resistance to interference are both highly desirable characteristics, especially with regard to those components of a wafer detection system which are located in the process chamber. Also, reducing the cost of wafer fabrication machines is a major concern in an industry where fabrication processes are machine intensive.\nIn attempts to solve these problems, the prior art employs sensors for detecting whether a wafer sits flat or inclined within a susceptor pocket. Many of the prior art detection systems employ optical sensors which are subject to interference from the intense and varying light produced by the process chamber heat lamps. Also, the prior art systems are unnecessarily complex and expensive. In addition, conventional detection system components are subjected to the harsh environment, which consequently further increases replacement costs as compared with a simpler system.\nAccordingly, a simpler, more accurate, and less expensive wafer placement detection system is needed."}
{"text": "1. Field of the Invention\nThis invention relates to a novel class of organosilicon compounds. This invention also relates to organosiloxane compositions that cure by a platinum group metal catalyzed hydrosilation reaction and contain these compound as adhesion promoting additives. The compounds are particularly useful for achieving adhesion to a variety of substrates at curing temperatures below about 100.degree. C.\n2. Background Information\nOne type of curable organosiloxane composition comprises an organopolysiloxane containing ethylenically unsaturated organic groups and an organosilicon compound containing silicon-bonded hydrogen atoms, curable in the presence of a catalyst which comprises a group VIII metal or a compound or complex thereof. Usually, the catalyst is a compound or complex of platinum. Unfortunately, the adhesion of these type of compositions is generally poor, limiting their usefulness in many applications.\nSeveral methods are known to improve the adhesion of organosiloxane compositions, however most generally have limitations. For example, primers can be used to improve the adhesion, but their use requires a separate processing step which increases process time and cost. Additionally, the primer usually requires a heat-activated cure to be effective. Adhesion additives, which are added directly to the composition to promote adhesion, have therefore received considerable attention.\nUseful adhesion additives typically contain at least two functional groups. U.S. Pat. No. 3,772,026 and U.S. Pat. No. 3,873,334 teach acyloxy functional silanes which additionally contain silicon bonded hydrogen atoms or alkenyl radicals, respectively. Although the acyloxy functionality is very reactive, it liberates corrosive acetic acid upon reaction with hydroxyl groups present at the surface of the substrate or with upon contact with water.\nU.S. Pat. No. 4,196,273 teaches use of unsaturated alkoxysilanes as adhesion additives, however, adhesion is not demonstrated at temperatures below 100.degree. C. and only adhesion to glass is demonstrated at 100.degree. C. U.S. Pat. No. 4,329,273 teaches use of partially hydrolyzed alkoxysilanes containing ethylenically unsaturated groups bonded to the silicon atom. While satisfactory curing of these compositions containing these adhesion additives was demonstrated at 100.degree. C., the compositions did not cure at ambient laboratory temperatures.\nU.S. Pat. No. 4,082,726 teaches using organosilicon compounds containing epoxy functional dialkoxysilyl groups and at least one alkenyl group or silicon-bonded hydrogen atom as adhesion additives. The exemplified compositions are cured at temperatures between 100.degree.-200.degree.C.\nU.S. Pat. No. 4,087,585 teaches physical blends of epoxy functional alkoxysilanes with silanol functional fluids containing alkenyl functionality. Cure was obtained at elevated temperatures. Multifunctional organosilanes containing ethylenic unsaturation, epoxy functionality, and alkoxy functionality are taught in U.S. Pat. No. 4,732,932. U.S. Pat. No. 5,106,933 teaches use of mixtures of alkoxysilanes.\nU.S. Pat. No. 4,659,851, U.S. Pat. No. 4,719,262, and U.S. Pat. No. 4,906,686 teach reaction products of 1) unsaturated alcohols and ethers derived from these alcohols and 2) alkoxysilanes. A shortcoming of these reaction products is that the resultant Sew-O-C bond between the unsaturated alcohol and alkoxysilane would be susceptible to hydrolysis to the extent that the adhesion may be decreased upon exposure to moisture.\nOrganosilicon compounds obtained from the reaction of polyhydric alcohols with compounds containing silicon-bonded hydrolyzable groups are known in the art. U.S. Pat. No. 2,349,338 teaches silicic acid esters of monohydric and polyhydric alcohols, including various glycol and glycerol silicate ester derivatives, are useful as corrosion prevention additives for heat transfer liquids. The corrosion prevention is afforded by the hydrolysis product of these compounds.\nU.S. Pat. No. 2,776,307 claims silicon-containing esters of the type (R.sup.1 O).sub.3 Si(XR.sup.2).sub.n OSi(OR.sup.3).sub.3, where X can be oxygen, and n has a value of 2, 3, or 4. The compounds are said to be useful as hydraulic fluids and other lubricant purposes.\nU.S. Pat. No. 3,992,429 teaches novel alkoxysilane cluster compounds containing sterically hindered alkyl groups to improve hydrolytic stability. Intended uses include heat transfer fluids, hydraulic fluids, brake fluids, transmission fluids, and the like.\nU.S. Pat. No. 2,630,446 teaches polymeric silicates prepared from tetraalkyl orthosilicates and polyhydric alcohols. The silicates are intended for use as plasticizers, lubricants, and hydraulic fluid additives.\nU.S. Pat. No. 4,072,655 teaches reaction products of polyhydric alcohols and silanes containing from 1 to 4 hydrolyzable groups per molecule. The reaction products are used as non-slump additives for moisture curable organosiloxane compositions.\nThe preparation of alkoxysilanes from glycols and glycol monoethers is taught by H. G. Emblem and K. Hargreaves in the Journal of Inorganic Nuclear Chemistry, Vol. 30 (3), p. 721, 1968).\nThe preparation of esters derived from tetraalkyl orthosilicates and polyhydroxy alcohols by a transesterification reaction is described in British Patent No. 450,875.\nU.S. Pat. No. 3,029,269 teaches compounds of the type (RO).sub.3 SiOCR'.sub.2 (CH.sub.2).sub.n CR'.sub.2 -OSi(OR).sub.3, where n is 0-16 and R' is H or a lower aliphatic radical, and R is a saturated aliphatic radical of 1-16 carbons. The compounds are suggested as lubricants and hydraulic fluids and in general show low pour point, good viscosity and viscosity index properties as well as hydrolytic stability.\nFinally, Applicants' copending application Ser. No. 08/240,598 filed on May 10, 1944, U.S. Pat. No. 5,424,384 describes adhesion promoting additives for curable organosiloxane compositions that are reaction products of a polyhydric saturated alcohol with a silane or bis-silylalkane. The organosilicon reactant contains at least three alkoxy or enoloxy groups on each silicon atom.\nOne objective of this invention is to provide a class of novel organosiloxanes that are particularly useful as adhesion promoting additives in organosiloxane compositions, particularly those that cure by a platinum group metal catalyzed hydrosilation reaction at temperatures below about 100.degree.C."}
{"text": "The present invention relates to a loading ramp that is interchangeable with the tailgate of a pickup requiring no alterations of the tailgate nor the tailgate mounting devices."}
{"text": "Media delivery has historically followed a broadcast type model, where users/consumers all receive the same programming. Thus, any effects, cross-fades or other blending between subsequent clips or program elements are performed upstream of the consuming device, prior to being sent over the broadcast channel(s). As is generally appreciated, the addition of these effects produces a high quality experience for the user, and also provides natural and enhanced transitions between program elements. These enhancements can significantly improve and enrich the listening experience, and can be changed or modified depending upon the “mood” of the channel, the sequence of songs or clips being played, as well as the audience type, time of day, and channel genre. Typically, elements that require cross-fading, blending or other signal processing of two or more elements require precise synchronization and simultaneous playback of the elements to be processed. Thus, although in the 1960s and 1970s DJs would try to mix songs in real time, by “cueing up” the next song and starting its turntable a bit before the currently being played song ended, with the advent of digital media it has become the norm to perform such processing on a playlist of multiple songs or clips prior to broadcasting it, storing its result at the media provider or broadcaster's servers, and then send it over the broadcast channel.\nWith the introduction of media compression and file based delivery, various types of media are commonly downloaded directly to a user's device, such as, for example, an iPod, digital media player, MP3 player, PC, tablet, cellular phone, smart phone, etc., and various hybrid devices or devices with equivalent functionalities, without the benefit of upstream processing between media elements. This leads to a less satisfactory user experience upon user consumption or playback. A user simply hears one song stop, then hears a brief pause, then hears the next song begin. There is no “awareness” by the media playing device as to what the sequence is, no optimizations as to which song most naturally follows another in the playlist, no sense of the “feel” “mood” or tempo of the playlist or any segment of it, and each sequence of media clips is, in general, unique to each user and how they organize their respective playlists.\nAdditionally, many consumer type devices, cell phones, smart phones, tablets, etc. do not have the capability to perform simultaneous decode and presentation of media and elements so that they can be cross-faded or processed as played back in real time. Such devices, for example cell phones, typically have a single hardware decoder per media type, so that any type of cross-fade in real time would also require additional software based decoding for the other elements, which (i) has negative impact on battery life, and (ii) would also require the precise synchronization of two or more decoders.\nWhat is needed in the art are systems and methods to implement and facilitate cross-fading, blends, interstitials and other effects/processing of two or more media elements on a downstream device for various purposes so as to enhance the elistenign experience, and, for example, replicate to the extent possible the sound and feel of broadcast programming.\nWhat is further needed in the art are methods to perform such processing involving two or more elements on a downstream device, where only a single hardware decoder is available or where other system constraints are operative."}
{"text": "1. Field of the Invention\nThe present invention relates to a face panel for a color image display apparatus, a panel for a color image display apparatus using the face panel, and a color image display apparatus on which the panel is mounted.\n2. Description of the Related Art\nIt is difficult for a conventional image display apparatus using a Cathode Ray Tube (CRT) called a Braun tube to cope with requirements for an enlargement of a display screen due to its weight, a technical problem for an increase of size of the CRT, and the like. Accordingly, image display apparatuses using a flat panel called Flat Panel Display (FPD) are predominantly used as image display apparatuses at present. Furthermore, liquid crystal displays, a plasma displays, and the like are mainly used as image display apparatuses using FPD.\nHowever, since the liquid crystal display and the plasma display are partially inferior to the conventional image display apparatus using a cathode ray tube (CRT) called a Braun tube in image quality and a response speed, they are required to further improve them.\nThere is a field emission display (FED) having the same operation principle as that of the CRT to improve the image quality and the response speed of the FPD.\nThere are characteristics of brightness of a displayed image as important characteristics of a color image display apparatus. The characteristics of the brightness of the displayed image are generally evaluated by luminance and a contrast ratio. The contrast ratio is a luminance ratio of white (maximum luminance) and black (minimum luminance). Since the black luminance is affected by the reflection of external light from an image screen, it is assumed below that the contrast ratio means a contrast ratio in an incident condition that external light incident on the image screen (bright place).\nFurthermore, it is assumed that the contrast ratio in the bright place is the contrast ratio which is measured in a condition that external light of 150 lux is irradiated to the image screen at an angle of 45°.\nThe brightness and the contrast ratio of an image, which is displayed on the image screen and visually recognized by a viewer, depend not only on the level of luminance of the displayed image but also on the surface brightness of the display screen itself. That is, the brightness and the contrast ratio of the displayed image visually recognized by the viewer are determined by the balance between the sum total of the reflected lights when no image is displayed on the image screen, the visually recognized brightness of a fluorescent screen itself, and the like and the luminance of the displayed image which emits light by the fluorescent screen.\nThen, the image display quality of a color image display apparatus can be improved by improving the brightness and the contrast ratio of the displayed image.\nWhen the contrast ratio of the FED is improved, it is considered that the technique of the CRT, which has the same operation principle as that of the FED, can be used.\nTo improve the contrast ratio of the CRT, Japanese Patent Application Laid-Open Nos. 64-7457 and H07-179711 disclose a technique of using a color filter, which has the same color as that of a phosphor, between a substrate and a phosphor layer, the color filter having a function for causing the light emitted from the phosphor to selectively pass therethrough. Japanese Patent Application Laid-Open No. H09-27284 discloses a manufacturing method and like of a color filter of a plasma display. Furthermore, Japanese Patent Application Laid-Open No. 2000-164145 discloses a technique for providing a filter on a surface of a substrate to improve a contrast ratio of a CRT.\nJapanese Patent Application Laid-Open No. 64-7457 discloses a color filter which is composed of pigment particles having an average particle diameter of 0.005 μm to 0.07 μm and which has a film thickness of 15 nm to 250 nm.\nFurthermore, Japanese Patent Application Laid-Open No. 09-27284 discloses a technique of not only interposing a color filter between a phosphor and the substrate but also forming an external filter to the surface of the substrate (outer surface) opposite to the surface thereof on which the phosphor is formed (inner surface).\nJapanese Patent Application Laid-Open No. 09-27284 discloses to form a color filter whose transmittance to a light having a wavelength within the range of the maximum spectrum wavelength of emitted lights ±20 nm is higher than the transmittance to a light having a wavelength of 400 to 650 nm other than the above range in the lights light emitted from a phosphor layer.\nThe external surface filter is a filter for supplementing the light absorption of the color filter and has a maximum absorption in a constant wavelength band other than the range of ±20 nm of the maximum spectrum wavelength of the light emitted by a three-color phosphor layer.\nJapanese Patent Application Laid-Open No. H07-179711 discloses a pigment-dispersed liquid composition containing an inorganic pigment of 0.5 to 50 parts by weight, at least one kind of a dispersing agent selected from a group composed of sodium salt, ammonium salt or amine salt of acrylic resin, acrylic copolymer resin, polycarboxylic acids, and an aromatic sulfonic acid formalin condensate, and water or a solvent compatible with water. Furthermore, the publication discloses a method of manufacturing a color filter using the pigment-dispersed liquid composition.\nAlthough the contrast ratio can be increased when the transmittance of a substrate to visible light is increased to about 50%, a halation is caused in a portion where a white display is made in contact with a black display. To prevent this phenomenon, Japanese Patent Application Laid-Open No. 2000-164145 discloses a technique for setting the transmittances of a substrate and a filter disposed to an outer surface of the substrate to 60 to 80%, respectively."}
{"text": "Electrophotography, which is a method for visualizing image information by forming an electrostatic latent image, is currently employed in various fields. The term “electrostatographic” is generally used interchangeably with the term “electrophotographic.” In general, electrophotography comprises the formation of an electrostatic latent image on a photoreceptor, followed by development of the image with a developer containing a toner, and subsequent transfer of the image onto a transfer material such as paper or a sheet, and fixing the image on the transfer material by utilizing heat, a solvent, pressure and/or the like to obtain a permanent image.\nDuring the development stage of electrophotography, as toner is magnetically attracted to the magnetic development roller, a negative charge build up is caused on the individual toner particles. This is called the triboelectric charge. Since both the toner and roller are charged negatively, the toner is repelled towards the positively charged areas of the drum to create the latent image. The negatively charged areas of the drum also repel toner, leaving only the image on the drum in dusted toner particles.\nIn electrostatographic reproducing apparatuses, including digital, image on image, and contact electrostatic printing apparatuses, a light image of an original to be copied is typically recorded in the form of an electrostatic latent image upon a photosensitive member and the latent image is subsequently rendered visible by the application of electroscopic thermoplastic resin particles and pigment particles, or toner. Electrophotographic imaging members may include photosensitive members (photoreceptors) which are commonly utilized in electrophotographic (xerographic) processes, in either a flexible belt or a rigid drum configuration. Other members may include flexible intermediate transfer belts that are seamless or seamed, and usually formed by cutting a rectangular sheet from a web, overlapping opposite ends, and welding the overlapped ends together to form a welded seam. These electrophotographic imaging members comprise a photoconductive layer comprising a single layer or composite layers.\nWhen the triboelectric charge of a toner is low, a number of issues occur in electrophotography using that toner. For example, there may be color spittings during development (e.g., toner particles accumulate at baffle and nips between the roll and edge of housing of the machine during printing and release them suddenly during printing marking undesirable large marks on prints), low developer life, and background (e.g., spots caused by individual particles that you see on the permanent image) on the final image. These issues can be magnified when a toner's triboelectric charge decreases with age. These issues are more prevalent with magenta developer, because triboelectric charge is generally lower than for other colors. External additives have been used to minimize these issues, but it has previously not been possible to increase the triboelectric charge of magenta developer to that of other colors. Previous attempts to sue an internal charge control agent have not been effective, because magenta toner is often itself a powerful charge control agent. Thus, the use of large quantities of external additives, the majority of the additives being negatively charging silica, are generally required to yield a functional magenta toner design."}
{"text": "As smart phones and other portable devices increasingly become ubiquitous, and data usage increases, macrocell base station devices and existing wireless infrastructure in turn require higher bandwidth capability in order to address the increased demand. To provide additional mobile bandwidth, small cell deployment is being pursued, with microcells and picocells providing coverage for much smaller areas than traditional macrocells.\nIn addition, most homes and businesses have grown to rely on broadband data access for services such as voice, video and Internet browsing, etc. Broadband access networks include satellite, 4G or 5G wireless, power line communication, fiber, cable, and telephone networks."}
{"text": "Various information processing systems have been developed as a result of the rapid changes which have taken place in the information industry in recent years. Methods of recording and apparatus compatible with new information processing systems have been developed and adopted. Thermal transfer recording methods use apparatus which is light and compact, has little noise, and has excellent operability and maintenance characteristics. Moreover, since thermal transfer also allow coloring to be achieved easily, these methods are widely used.\nThermal transfer recording methods can be broadly classified into two types, namely mass transfer types and dye transfer types. The latter case relates to a recording method in which a thermal dye transfer donating material (hereinbelow, \"dye-donor\") is constucted of a substrate with a dye layer containing dyes having heat transferability. The material is brought into contact with a thermal dye transfer receiving material (hereinbelow, \"dye receptor\"). The dye donor material is selectively heated with a thermal printing head provided with a plurality of juxtaposed heat-generating resistors. The heating is in response to an information signal defining a pattern or image. Dye from the selectively heated regions of the dye donor is transferred to the dye receptor and forms a pattern thereon. The shape and the density of the patern forms an image in accordance with the pattern and the intensity of heat applied to the dye-donor.\nA dye receptor usually comprises a support coated with a dye receiving layer. The dye coming from the dye donor can thermally and properly diffuse into that layer. An intermediate layer, useful as cushioning layer, porous layer or dye diffusion preventing layer, may be provided between the support and the receiving layer.\nThe dye donor may be a monochrome dye layer or it may comprise a sequence of different colored and discrete areas of, for example, cyan, magenta, yellow, and optionally black hue. When a dye-donor containing said sequenced two, three or more primary color areas is used, a multicolor image can be obtained by sequentially performing the dye transfer process steps for each color. The dye receptors of the prior art are commonly manufactured by coating organic solvent solutions of polymers and other ingredients, involving expensive, polluting and hazardous processes. To reduce risks of fire, explosions and other accidents, special precautions and expensive manufacturing apparatus are needed in handling the organic solvent solutions used in that type of manufacture.\nThe image fastness given by the prior art dye receptors is quite limited and still not competitive with conventional photographic image fastness.\nTo bypass the use of organic solvents, JP Patent Appls. 57/137,191 or 60/038,192 claims dye receptors obtained by coating a blend of polyesters or vinylic latices that however still give the disadvantage of limited image fastness, including significant photofading.\nEuropean Appl. 363,989 describes dye receptors based on water soluble polymers in which polymeric dye accepting compounds are dispersed, and wherein said water soluble polymers are hardened by a hardening agent.\nSimilarly, JP Patent Appl. 02/025,393 describes dye receptors based primarily on polymer solutions as a primary binder and vinyl styrene or ethylvinylacrylate particles as a secondary ingredient.\nEP 351,075 is another prior art example of aqueous dye receptors, using a silica dispersion and a melamine and formaldehyde condensation resin. In EP 300,505 a polyolefin latex is used to coat a receptor underlayer. The dye receiving layer is obtained by coating an organic solvent solution of polymer.\nIn JP Patent Appl. 61/266,296, aqueous receptors are obtained by using aqueous solutions of water soluble polymers such as polyvinyl alcohol or substituted celluloses as a binder for porous and non-porous fillers.\nIn JP Patent Appl. 63/315,283, aqueous solutions of polyvinyl alcohol and/or other water soluble resins are used as receptor binders. In EP 364,900 a polyester receptor layer is obtained by polycondensation of polyfunctional acids and alcohols and curing of the aqueous coated solution of reactants to crosslink them.\nIn DE 3,934,014 copolymers of styrene and acrylic compounds are used as latices for obtaining the underlayer. The dye receiving layer is coated over the latex underlayer.\nJP 02/122,992 discloses a receiving layer comprising an aqueous solution or dispersion of polymeric resin in combination with silica particle and modified silicone oil, the layer having improved antisticking properties.\nJP 01/038,277 discloses a composition for a receiving layer obtained from an aqueous dispersion of modified polyester containing hydrophilic groups.\nIn JP 01/004,391 aqueous latices with a Tg&gt;50.degree. C. are involved in the preparation of dye receptors in combination with colloidal silica.\nJP 63/011,392 discloses an oil solution of resin dispersed in water and then coated.\nIn JP 62/238,790 a solution or dispersion of polyester having solubilizing groups is combined with a water solution or dispersion of resins and of crosslinking compounds to increase the adhesion of the receiving layer.\nIn JP 62/146,693 a latex is coated as an underlayer (cushioning layer) on which the receiving layer is coated.\nAccordingly, there is at present continuous work to obtain aqueous dye receptors with improved qualities which reduce the above mentioned problems."}
{"text": "In a mobile communication system, if a building is built in the vicinity of a base station, or if the installation status of a base station in the vicinity of the base station changes, the radio communication environment according to the base station changes. Therefore, conventionally, the operators make use of a measurement vehicle in which measurement equipment is loaded, and perform a drive test by measuring a state of a signal received from the network (base station), and then collecting the measurement data.\nSuch a measurement and collection process can, for example, contribute to the optimization of coverage of the base station, but faces the problem of too many man-hours and high cost. Thus, according to 3GPP (3rd Generation Partnership Project), i.e., a standardization project for mobile communication systems, a radio terminal belonging to the user is used to proceed with the specification planning of MDT (Minimization of Drive Test) for automation of the measurement and collection.\nA type of MDT is the immediate report type MDT (referred to as “Immediate MDT”). According to Immediate MDT, a radio terminal in the connected state measures a reception signal state, and then reports the measurement results and location information during measurement to the network. The connected state refers to the state when the radio terminal is performing communications.\nFurthermore, it is stipulated in the 3GPP standards that when a radio terminal in the connected state detects a connection failure with the network, a radio link failure (RLF) report about the connection failure is sent from the radio terminal to the network (see 3GPP TS 36.300 v10.4.0). A connection failure refers to RLF or a handover failure. Furthermore, it is stipulated that in Immediate MDT, the RLF report is extended, and the location information for identifying the location where the connection failure occurred is included in the RLF report (see 3GPP TS 37.320 v10.1.0).\nBased on the Immediate MDT report and the RLF report received from the radio terminal, the network performs parameter adjustment (settings change) for coverage optimization."}
{"text": "The present invention relates to a network relaying apparatus and a network relaying method, or in particular to a network relaying apparatus including a router of a computer network system which is capable of searching at high speed for a destination of a packet input and a network relaying search method.\nGenerally, in a network system, a network relaying apparatus such as a router or a bridge is used for connecting a plurality of networks. The router checks the destination address of a packet received from a network or a subnet connected, determines the destination of the packet, and transfers the packet to a network or a subnet which is connected with the destination router or host.\nFIG. 13 is a diagram showing a configuration of a conventional network relaying apparatus. In FIG. 13, a router 100 includes a routing manager (RM) 110, router buses 120, network interfaces (NIF) 130 and ports 140. Each port 140 is connected to an appropriate network 150.\nEach network interface 130 receives a packet from a network connected to the port 140, and transmits the received packet through the router bus 120 to the routing manager 110. The routing manager 110 includes a routing table for holding the routing information, and using this routing information, determines the network 150 of the destination from the address of the packet received, and transmits the packet to the network interface 130 of the port 140 connected to the network 150. The network interface 130 that has received the packet from the routing manager 110 sends out the packet to the destination network 150. The routing manager 110 updates and maintains the routing information held in the routing table based on the header information of the packet received, and has the function of overall management of the router 100.\nAn explanation will be given of the route search process for searching for a port outputting the next address to which the packet is to be transferred upon receipt of the packet and outputting the packet. Normally, the route search uses a route search table (routing table) prepared from the component definition information and the information obtained by exchange between the routers. The routing table is for searching the information (next hop information) as to the output port, the next hop address and whether the network is directly connected or not with a set of the network address and the network mask length as a key.\nAs another conventional system, JP-A-05-199230 (U.S. Pat. Ser. No. 5,434,863) discloses an internet-work system and a communication network system which can flexibly meet the size requirement of the network without adversely affecting the high-speed routing process. In these systems, a router manager and a plurality of routing accelerator modules are coupled to each other with a high-speed bus Also, each routing accelerator is connected with a plurality of independent communication ports. In these conventional systems, a plurality of the routing accelerators makes possible a high-speed routing and by adding the routing accelerators, the requirement for increasing the network size can be easily met.\nThe conventional router, however, cannot meet the requirement of the high speed lines such as the high-speed LAN (local area network) and the wide band ISDN (Integrated Services Digital Network) and ATM (asynchronous transfer mode) that have recently found applications. Also, the conventional router with only one routing means has the disadvantage that the number of ports and the communication traffic that can be supported are limited. It is therefore difficult to expand the configuration of the port menu of the router to a large size smoothly or to improve the performance in keeping with the port traffic volume.\nAlso, with the increase of internet users, the number of flows that the router is required to detect is on the increase. Therefore, it is necessary to set a multiplicity of flow conditions in the router. The increased traffic and the increased line speed of the internet, on the other hand, requires a shorter processing time per packet in the router. Even in the case where the number of set flow conditions is increased, therefore, the QoS (quality of service) control and the filter operation must be carried out at a high speed on the part of the router.\nIn setting flow conditions, on the other hand, a great variety of flow conditions set as desired by the router manager must be flexibly handled. The prior art fails to take this point into account.\nIn view of the above-mentioned points, an object of the present invention is to provide a network relaying apparatus and method for routing packets at high speed while assuring a high communication quality of service (QoS), a high reliability and security.\nAnother object of the invention is to provide a network relaying apparatus and method in which the flow conditions including the information for identifying users, the protocol information and the priority information can be set in great amounts, and in keeping with the increase in the line speed and the flow conditions, the flow can be detected at high speed so that the control operation for the communication quality including QoS control and filter can be realized.\nStill another object of the invention is to provide a network relaying apparatus and method in which the control operation flexibly meeting a great variety of flow conditions including the priority control, discard control and band control can be performed at high speed by improving the description of the flow conditions and the combination of the information including the source and the transfer destination.\nAccording to this invention, as described above, there are provided a network relaying apparatus and method for routing packets at high speed while at the same time assuring a high communication quality (QoS), a high reliability and tight security.\nAlso, according to this invention, the flow conditions including the information for identifying the users, the protocol information and the priority information can be set in great amounts in keeping with the increase in line speed and flow conditions, and the flow can be detected at high speed for realizing a high-speed QoS control and filtering. Further, the high-speed control operation flexibly meeting a great variety of flow conditions including the priority control, the discard control and the band control is made possible by improving the descriptiveness of the flow conditions and the combination of the information including the source and the transfer destination. Further, according to this invention, jobs of different categories (such as basic jobs and information jobs) can be combined into a single network.\nAccording to one aspect of the invention, there is provided a network relaying apparatus comprising:\nat least a network interface connected with at least a network;\nat least a routing processor including a packet buffer for storing input packets and a flow search table set separately for each of the input or output line number with action information corresponding to the information including the packet source and the packet transfer destination as an entry;\na routing manager for managing the internal components of the system; and\na connector for connecting the routing manager and each of a plurality of the routing processors;\nwherein the network interface outputs the input packet from the network to the routing processor; and\nwherein the routing processor includes means for storing the input packet from the network interface in a buffer memory, means for searching the transfer destination of the input packet stored in the packet buffer based on the stored header information, means for searching and reading only the entry corresponding to the input or output line number of the packet by referring to the flow search table, means for determining whether the information including the packet source and the packet transfer destination are coincident with the reference conditions in the entry read out, means for determining, in the case of coincidence, the control operation for the communication quality including the order of priority of packet transfer and the possibility of transfer in accordance with the action information in the entry, and means for outputting the input packet stored in the packet buffer and the output packet produced according to the header information, to the connector and the network interface.\nAccording to another aspect of the invention, there is provided a network relaying method for outputting the input packet input from a network to a transfer destination in a network relaying apparatus comprising at least a network interface connected to at least a network, at least a routing processor for routing the packet input from the network interface, a routing manager for managing the internal components of the system, and a connector for connecting the routing manager and each of a plurality of the routing processors;\nwherein the routing processor includes:\nmeans for setting the action information corresponding to the information including the packet source and the packet transfer destination as an entry separately for each input or output line number in a flow search table;\nmeans for storing the input packet in a buffer memory;\nmeans for searching for a transfer destination of the input packet stored in the packet buffer based on the stored header information;\nmeans for searching and reading only the entry corresponding to the input or output line number of the packet by referring to the flow search table;\nmeans for determining whether the information including the packet source and the packet transfer destination coincides with the reference conditions in the read entry;\nmeans for determining, in the case of a coincidence, the control operation for the communication quality including the priority of packet transfer or the possibility of transfer based on the action information in the entry; and\nmeans for outputting the input packet stored in the packet buffer and the output packet produced by the header information to the connector or the network interface."}
{"text": "This invention relates to the residue number system (RNS) and more particularly to improved apparatus for converting residue number signals to binary number signals which are readily converted to analog representations of decimal number signals.\nIt is becoming increasingly popular to design electrical equipment such as filters, CODECS and echo cancellers to perform necessary processing functions in a digital manner in order to obtain the advantages of high reliability and reproducibility of results. This requires converting an analog input signal to a digital signal, performing the necessary processing operations on the digital signal in a digital manner, and then translating the processed digital signal back into an associated analog output signal. In a digital filtering operation, for example, a series of multiplications and a long summation of product terms may be required. Since there are carries involved with basic arithmetic operations such as multiplication and addition in weighted number systems such as the decimal and binary (base 10 and base 2) number systems, time is required to accommodate all of the necessary carry operations. By way of example, there is a single carry from the units column to the 10's column when adding the numbers 9 and 5 to get the number 14 in the base 10 number system. This also applies to multiplication. And where 2 each 16 digit numbers are multiplied together to obtain a 32 digit number, it may be necessary to propagate a carry from the least significant digit in the units column all the way up to the last digit that corresponds to the most significant digit. This can be very time consuming. One technique for overcoming this time limitation that is called pipelining involves generation, in the same time, of different future product (or addition) terms corresponding to different bits and delaying them accordingly, so that all bits of a result will appear on outputs at the same time. This normally requires considerable hardware which can be costly in terms of both economics and real estate in an integrated circuit structure. Another technique for overcoming this time limitation is to perform such operations as multiplication and addition in a non-weighted number system which is called the residue number system (RNS). A residue signal processor has advantage over a conventional binary-digital processor in that results are exact, a smaller amount of hardware is required, shallower pipelining is possible, operation time may be decreased, and no carries are involved.\nDuring recent years there has been considerable interest in and development of techniques for performing digital computing and digital signal processing with hardware based on the RNS. This interest stems at least partially from the fact that operations can be performed quickly in the RNS because of the separability of operations on each RNS digit and the elimination of carries. Residue signal processor means must sample analog input signals and convert the samples to residue signals, perform the necessary processing on the residue signals, and convert the resultant residue signals back to analog output signals. Although techniques are readily available for efficiently converting analog/decimal or binary/digital signals to residue number signals, it is much more difficult to convert residue number signals back into analog signals. Since the RNS is not a weighted number system, its direct conversion to decimal/analog is not practical. One approach is to convert residue number signals to digital/binary signals and use one of the well known digital to analog conversion schemes to obtain the requisite analog/decimal signals.\nResidue to analog conversion is described in the articles:\n\"Techniques for Residue-to-Analog Conversion for High Data Rate Digital Filtering\" by W. K. Jenkins, 1978 IEEE International Conference on Acoustics, Speech and Signal Processing, pages 804-807;\n\"An Efficient Residue-to-Decimal Converter\" by F. J. Taylor and A. S. Ramnarayanan, IEEE Transactions on Circuits and Systems, Vol. 28, No. 12, December 1981, pages 1164-1169;\n\"An Improved Residue Number System Digital-to-Analog Converter\" by M. A. Soderstrand, et al, IEEE Transactions on Circuits and Systems, Vol. 30, No. 12, December 1983, pages 903-907; and\n\"Techniques for Residue-to-Analog Conversion for Residue-Encoded Digital Filters\" by W. K. Jenkins, IEEE Transactions on Circuits and Systems, Vol. 25, No. 7, July 1978, pages 555-562.\nPrior art residue to binary/digital converters generally require considerable hardware and/or have relatively slow execution times and/or require considerable surface area, which is undesirable in an integrated circuit structure or implementation. Also, many of the prior art residue to analog converters require memories that store look-up tables and/or modulo adders, both of which are generally hardware intensive, slow in terms of overall execution time and/or are magnitude limited in that memory size dictates the maximum number of bits and thus the largest number that can be handled. The execution time of RNS to binary converters which require memory is normally dictated by the access time of the memory, (i.e., the time that it takes to obtain information from a memory location) which may be for CMOS ROMS in the order of 100 nanoseconds. Memory systems having much shorter execution times are very expensive and are even more magnitude limited (i.e., memories with short execution times normally have fewer bytes of memory). Stated differently, memories having larger storage capacity normally have longer execution times and require considerably more real estate, which is a disadvantage in an integrated circuit structure. This can be a significant limitation since the advantages of the RNS increase with the number of bits of memory that are available. By way of example, multiplication of 2 each n-bit numbers is n times faster in the RNS than in the binary number system without pipelining. Additionally, modulo adders are normally more complex than and require more surface area in integrated circuit implementations than do conventional binary adders.\nAn object of this invention is the provision of improved method and apparatus for converting residue number signals to binary number signals."}
{"text": "Night vision systems are used in a wide variety of military, industrial and residential applications to enable sight in a dark environment. For example, night vision systems are utilized by military aviators during nighttime flights. Security cameras use night vision systems to monitor dark areas and medical instruments use night vision systems to alleviate conditions such as retinitis pigmentosis (night blindness).\nImage intensifier devices are employed in night vision systems to convert a dark environment to an environment perceivable by a viewer. More specifically, the image intensifier device within the night vision system collects tiny amounts of light in a dark environment, including the lower portion of the infrared light spectrum present in the environment, which may be imperceptible to the human eye. The device amplifies the light so that the human eye can perceive the image. The light output from the image intensifier device can either be supplied to a camera, external monitor or directly to the eyes of a viewer. The image intensifier device is commonly employed in vision goggles that are worn on a user's head for transmission of the light output directly to the viewer. Accordingly, since the goggles are worn on the head, they are desirably compact and light weight for purposes of comfort and usability.\nImage intensifier devices include three basic components mounted within a housing, i.e. a photocathode (commonly called a cathode), a microchannel plate (MCP), and a phosphor screen (commonly called a screen, fiber-optic or anode). The photocathode detects a light image and converts the light image into a corresponding electron pattern. The MCP amplifies the electron pattern and the phosphor screen transforms the amplified electron pattern back to an enhanced light image.\nReferring to FIG. 1, a current state of the art Generation III (GEN III) image intensifier tube 10 is shown. Examples of the use of such a GEN III image intensifier tube in the prior art are exemplified in U.S. Pat. No. 5,029,963 to Naselli, et al., entitled REPLACEMENT DEVICE FOR A DRIVER'S VIEWER and U.S. Pat. No. 5,084,780 to Phillips, entitled TELESCOPIC SIGHT FOR DAYLIGHT VIEWING. The GEN III image intensifier tube 10 shown, and in both cited references, is of the type currently manufactured by ITT Corporation, the assignee herein. In intensifier tube 10 shown in FIG. 1, infrared energy impinges upon photocathode 12. The photocathode 12 is comprised of glass faceplate 14 coated on one side with antireflection layer 16, a gallium aluminum arsenide (GaAlAs) window layer 17 and gallium arsenide (GaAs) active layer 18. Infrared energy is absorbed in GaAs active layer 18, thereby resulting in the generation of electron/hole pairs. The produced electrons are then emitted into vacuum housing 22 through a negative electron affinity (NEA) coating 20 present on GaAs active layer 18.\nA microchannel plate (MCP) 24 is positioned within vacuum housing 22, adjacent NEA coating 20 of photocathode 12. Conventionally, MCP 24 is made of glass having a conductive input surface 26 and a conductive output surface 28. Once electrons exit photocathode 12, the electrons are accelerated toward input surface 26 of MCP 24 by a difference in potential between input surface 26 and photocathode 12 of approximately 300 to 900 volts. As the electrons bombard input surface 26 of MCP 24, secondary electrons are generated within MCP 24. The MCP 24 may generate several hundred electrons for each electron entering input surface 26. The MCP 24 is subjected to a difference in potential between input surface 26 and output surface 28, which is typically about 1100 volts, whereby the potential difference enables electron multiplication.\nAs the multiplied electrons exit MCP 24, the electrons are accelerated through vacuum housing 22 toward phosphor screen 30 by a difference in potential between phosphor screen 30 and output surface 28 of approximately 4200 volts. As is the electrons impinge upon phosphor screen 30, many photons are produced per electron. The photons create the output image for image intensifier tube 10 on the output surface of optical inverter element 31.\nFIG. 2 is a schematic representation of image intensifier tube 41. The tube includes photocathode 54, microchannel plate (MCP) 53 and imaging sensor 56. Imaging sensor 56 can be any type of solid-state imaging sensor, such as a CCD device, or a CMOS imaging sensor.\nPhotocathode 54 can be, but is not limited to, a material such as GaAs, Bialkali, InGaAs, and the like. Photocathode 54 includes input side 54a and output side 54b. MCP 53 has a plurality of channels 52 formed between an input surface and an output surface.\nAn electric biasing circuit 44 provides a biasing current to image intensifier tube 41. Electric biasing circuit 44 includes a first electrical connection 42 and a second electrical connection 43. First electrical connection 42 provides a biasing voltage between photocathode 54 and MCP 53. Second electrical connection 43 applies a biasing voltage between MCP 53 and imaging sensor 56. In this configuration, photocathode 54, MCP 53, and imaging sensor 56 are maintained in a vacuum body or envelope 61 as a single unit, in close physical proximity to each other.\nStill referring to FIG. 2, in operation, light 58, 59 from an image 57 enters image intensifier tube 41 through input side 54a of photocathode 54. Photocathode 54 changes the entering light into electrons 48, which are output from output side 54b of photocathode 54. Electrons 48 exiting photocathode 54 enter channels 52 of MCP 53. Secondary electrons are generated within the plurality of channels 52 of MCP 53. The MCP 53 may generate several hundred electrons in each of channels 52 for each electron entering through the input surface. Thus, the number of electrons 47 exiting channels 52 is significantly greater than the number of electrons 48 that entered channels 52. The intensified number of electrons 47 exit channels 52 and strike the electron receiving surface of imaging device 56. The imaging device transforms the electrons into a light image which may be stored in memory or viewed on display 46."}
{"text": "This invention relates generally to electrical connectors and more particularly to a molded connector having a construction which facilitates the assembly of the electrical connections for a float controlled pump such as a sump pump.\nElectric sump pumps are commonly equipped with a float having a mercury switch or mechanical switch that automatically turns the pump motor on when the liquid in the sump rises to a preselected level. In the past, the electrical lead wires from the motor, the power cord and the float switch cord have all been extended into an electrical compartment on the pump housing, and the proper wires must then be connected with one another by hand in the electrical compartment; another arrangement that has been used involves the provision of a separate cord with the float switch attached and a separate plug with \"piggy-back\" connector for the purpose of attaching the pump power cord. As can easily be appreciated, this assembly technique requires considerable time and a reasonable amount of skill, and it is also possible for the electrical connections in the box to be faulty.\nThe present invention is directed to an electrical connector which is constructed in a unique manner in order to facilitate assembly of the electrical wiring, reduce the assembly time, and virtually eliminate faulty electrical connections in a float controlled electric pump. In accordance with the invention, a tee shaped electrical connector is molded in a single integral piece with the power cord molded to one leg of the connector body and the float switch cord molded to another leg. The final leg of the connector body has a plug head equipped with electrical contacts to which lead wires from the power cord and the float switch cord extend. The pump casing is provided with a socket into which electrical terminals for the motor leads project. The plug head may be inserted into the sockets so that the electrical contacts mate with the terminals, and the connector body may be fastened in place by two screws.\nDue to this construction, all that is required for the electrical connections to be completed during the assembly process is to plug the connector into the socket and secure it with the screws. Because the connector body is prewired, the electrical circuit that results following assembly is properly arranged so that the float switch controls the on/off condition of the pump motor in the intended manner."}
{"text": "The invention relates to compositions suitable for use as adhesives, especially for medical products. This invention also relates to medical devices that include portions bonded together by the adhesives.\nAdhesives can be used in the medical field in a variety of ways. Medical adhesives can be used, for example, in manufacturing and/or implantation of medical devices. In surgical settings, medical adhesives provide an alternative to sutures, staples and the like for closing wounds in soft tissue. Certain tissues, such as nerves and particular vital organs, are too delicate for suturing or stapling, so the use of medical adhesives may be a viable repair option. In manufacturing settings, medical adhesives can be used to adhere one or more components of a medical device together.\nGenerally, the use of an adhesive for the manufacturing and/or implantation of medical devices or to repair wounds in soft tissue can be desirable due to potential sealing properties and uniform stress distribution. Adhesives provide a means for attaching substrates in the medical devices. The use of an adhesive may form an adhesive bond over the length of an attachment site. Adhesives tend to distribute the stress exerted on the attachment site over the length of the attachment site.\nGenerally, medical adhesives can be classified according to whether they include synthetic polymers, natural (biological) polymers or both. A variety of synthetic urethane based polymers have been developed as medical glues. Natural surgical adhesives generally are based on proteins. For example, fibrin glues include the protein fibrinogen. Fibrinogen is used in natural wound healing mechanisms in humans and other mammals. Synthetic adhesives have the disadvantage of being potentially toxic. On the other hand, biological/natural adhesives generally have relatively low binding (cohesive) strengths and rapid degradation times.\nIn a first aspect, the invention pertains to an adhesive composition including an epoxyamine and a curing agent. The epoxyamine is at a concentration of between about 25 percent by weight and about 90 percent by weight in the adhesive composition.\nIn a further aspect, the invention pertains to a prosthesis including a first substrate. The first substrate is associated with an adhesive composition that includes an epoxyamine. The concentration of the epoxyamine in the adhesive composition is between about 25percent by weight and about 90 percent by weight.\nIn another aspect, the invention pertains to method of securing a first substrate to a second substrate. The method includes contacting a first substrate including an adhesive composition with a second substrate to form an adhesive bond between the first and the second substrate. The adhesive composition includes an epoxyamine.\nIn a further aspect, the invention pertains to a method of manufacturing a heart valve prosthesis. The method includes associating a first valve component with an adhesive composition and adhering the first valve component with a second valve component to form an adhesive bond between the first valve component and the second valve component. The adhesive composition includes an epoxyamine.\nIn another aspect, the invention pertains to a prosthesis including at least two substrates. The substrates are adhered together by an adhesive bond. The adhesive bond is formed when an epoxyamine adhesive between the substrates is cured."}
{"text": "Signal transduction is the general process by which cells respond to extracellular signals. Signal transduction across the plasma membrane begins with the binding of a signal molecule, e.g., a hormone, neurotransmitter, or growth factor, to a cell membrane receptor. The receptor, thus activated, triggers an intracellular biochemical cascade that ends with the activation of an intracellular target molecule, such as a transcription factor. This process of signal transduction regulates all types of cell functions including cell proliferation, differentiation, and gene transcription.\nBiological membranes surround organelles, vesicles, and the cell itself. Membranes are highly selective permeability barriers made up of lipid bilayer sheets composed of phosphoglycerides, fatty acids, cholesterol, phospholipids, glycolipids, proteoglycans, and proteins. Membranes contain ion pumps, ion channels, and specific receptors for external stimuli which transmit biochemical signals across the membranes. These membranes also contain second messenger proteins which interact with these pumps, channels, and receptors to amplify and regulate transmission of these signals."}
{"text": "The present invention relates generally to an exposure method and apparatus used to expose an object, such as a single crystal substrate for a semiconductor wafer and a glass plate for a liquid crystal display (“LCD”)\nReduction projection exposure apparatuses have been conventionally employed which use a projection optical system for projecting a circuit pattern formed on a mask (reticle) onto a wafer, etc. and for transferring the circuit pattern, in manufacturing such devices as semiconductor devices, LCD devices and thin-film magnetic heads in the photolithography technology.\nAlong with finer and more highly integrated ICs, the projection exposure apparatus has been required to expose a circuit pattern on a reticle onto a wafer with high resolving power. The minimum critical dimension to be transferred by the projection exposure apparatus (resolution) is proportionate to a wavelength of light used for exposure, and inversely proportionate to the numerical aperture (“NA”) of the projection optical system. The shorter the wavelength is, the better the resolution is. Accordingly, the recent light sources have been in transition from an ultra-high pressure mercury lamp (g-line with a wavelength of approximately 436 nm) and i-line with a wavelength of approximately 365 nm) to KrF excimer laser (with a wavelength of approximately 248 nm) and ArF excimer laser (with a wavelength of approximately 193 nm). Practical use of F2 excimer laser (with a wavelength of 157 nm) has been promoted.\nA step-and-repeat exposure apparatus (also referred to as a “stepper”) for entirely projecting and exposing an approximately square exposure area onto a wafer with a reduced exposure area has been replaced mainly with a step-and-scan exposure apparatus (also referred to as a “scanner”) for accurately exposing a wide screen of exposure area through a rectangular slit with relatively and quickly scanning the reticle and the wafer.\nIn exposure, the scanner uses a surface-position detector having an oblique light projection system to measure a surface position at a certain position on the wafer before the exposure slit area moves to the certain position on the wafer, and accords the wafer surface with an optimal exposure image-surface position when exposing the certain position, thereby reducing influence of the flatness of the wafer. In particular, there are plural measurement points in longitudinal direction of the exposure slit, i.e., a direction orthogonal to the scan direction, at front and back stages to the exposure slit area to measure an inclination or tilt of the surface as well as a height or focus of the wafer surface position. In general, the scan exposure proceeds in both directions from the upper stage and from the back stage (see, for example, Japanese Patent Application Publication No. 9-45609 corresponding to U.S. Pat. No. 5,750,294).\nJapanese Patent Application Publication No. 6-260391 (corresponding to U.S. Pat. No. 5,448,332) proposes, as a method for measuring a surface position on a wafer in a scanner and for correction the same, an arrangement of plural measurement points on a pre-scan area other than the exposure area to measure the focus and tilt in scan and non-scan directions. Japanese Laid-Open Patent Application No. 6-283403 (corresponding to U.S. Pat. No. 5,448,332) proposes as a method for measuring the focus and tilt in the scan and non-scan directions and for driving and correcting the same, by arranging plural measurement points in the exposure area.\nA description will be given of these proposals with reference to FIGS. 18 and 19. Here, FIG. 18 is a schematic sectional view of focus and tilt measurement points FP1 to FP3 on the wafer 1000. FIG. 19 is a schematic sectional view showing the wafer 1000 that has been driven to an optimal exposure image-surface position based on the measurement results. Referring to FIG. 18, the focus and tilt are sequentially measured at the measurement points FP1 to FP3 on the wafer 1000. A pre-scan plane PMP is calculated based on the measurement results from the measurement points FP1 to FP3, and the orientation of the wafer is driven and adjusted to the best focus plane BFP in moving the wafer 1000 to the exposure are an exposure slit 2000, as shown in FIG. 19.\nHowever, the recent increasingly shortened wavelength of the exposure light and the higher NA of the projection optical system have required an extremely small depth of focus (“DOF”) and a stricter accuracy with which the wafer surface to be exposed is aligned to the best focus position BFP or so-called focus accuracy.\nIn particular, they have also required stricter measurement and precise correction of the tilt of the wafer surface in the scan direction or width direction of the exposure slit. A wafer having an insufficiently flat surface has disadvantageous focus detection accuracy. For example, when the exposure apparatus has a DOF with 0.4 μm, the flatness of the wafer requires several nanometer order, for example, the flatness of the wafer needs 0.08 μm where it is one-fifth as long as the DOF, or 0.04 μm where it is one-tenth as long as the DOF. In addition, while a surface-position detector having the oblique light projection system measures the wafer's surface position before the area hangs over the exposure slit, the measurement timing is discrete and no information is available or considered about the wafer's flatness between two timings. As a result, there is no information available between timings of the flatness of the wafer.\nFor example, this measurement timing is at an interval of 2 mm on the wafer 1000 in the scan direction as shown in FIG. 20. Then, the wafer 1000 has such an insufficient flatness due to lack of the information for a distance of 2 mm, e.g., between points P1 to P2 in FIG. 20 that the front position may offset by Δ from the pre-scan plane PMP calculated by the measurement at the interval of 2 mm. Here, FIG. 20 is a schematic sectional view showing an offset of flatness between the pre-scan plane PMP and the wafer 1000.\nIn exposure, the pre-scan plane PMP is adjusted to the best focus plane BFP, and the exposure in FIG. 20 needs a shift by the amount of Δ. This shift occurs in a direction orthogonal to the scan direction as well as the scan direction. This is due to an arrangement of measurement points of the above oblique light projection system, rather than the measurement timing.\nThe finer measurement timing in the scan direction and the increased number of measurement points in the oblique light projection system would reduce an offset error, but might disadvantageously lower the throughput due to the deteriorated scan speed in exposure time, increase measurement time, rise cost together with the complicated apparatus structure, and grows likelihood of troubles."}
{"text": "1. Field of the Invention\nThe present invention concerns a cooler for electrical and/ or electronic components, more particularly a cooler used in a testing system for semiconductors.\n2. Discussion of the Background Art\nThe lifetime and proper functioning of electric and/or electronic components, in particular microprocessors, depends heavily on their thermal strain. In order to increase the lifetime and to ensure optimum functioning, coolers of the above-mentioned type are used for cooling such thermally strained components.\nCoolers for the most varied applications and in different embodiments are known from, among others, JP-A-07038025, DE-A-19600166, JP-A-2000323636, DE-A-19800166, and DE-U-9201158."}
{"text": "The present disclosure generally relates to a storage system that provides easier access to multiple portions of a storage rack.\nStorage systems are used in warehouses, department stores, and storage facilities to store products thereon. The storage systems containing a plurality of storage racks may hold large amounts and a variety of products that consumers or pickers pull from. Some storage systems are designed in a manner to provide support with insufficient accessibility, thereby requiring additional resources and time for workers to pull all of the products they require to fill a particular order. For example, some storage systems include a horizontal bar that extends to another storage system adjacent thereto; the horizontal bar provides support for both storage system structures. However, the horizontal bar obstructs the pathway for the consumer or picker to access products on an opposite side of the horizontal bar."}
{"text": "This invention comprises a new and distinct variety of Orchid named Dendrobium Stardust xe2x80x98Firebirdxe2x80x99, a hybrid of the Genus Dendrobium (Dendrobium unicomxc3x97Dendrobium ukon), referred to herein by its grex and varietal name Stardust xe2x80x98Firebirdxe2x80x99. Stardust xe2x80x98Firebirdxe2x80x99 is a distinct variety selected from progeny of a cross between the unpatented Dendrobium unicum and Dendrobium unkon. The same cross has produced other Stardust family members (see below).\nThe chronology of events leading to selection of Stardust xe2x80x98Firebirdxe2x80x99 is as follows: Dendrobium unicum and Dendrobium unkon were crossed in 1980 and seedlings were placed in flasks in October of that year. These seedlings were grown to maturity, flowered, and the Grex epithet Stardust was registered with The Royal Horticultural Society, The International Registration Authority for Orchid Hybrids, in 1986 by N. Ashi. From this group of seedlings a plant with desirable characteristics (i.e., good flowers, compact growth habit, and fast growth), was selected and given the varietal epithet xe2x80x98Firebirdxe2x80x99.\nThis plant was then reproduced asexually using conventional meristem culture techniques for tissue culture of orchids (Plants from Test Tubes, Lydiane Kyte and John Kleyn, eds., Timber Press (1996)). Asexual reproduction took place at 60 Aza Hanto, Oaza Ishihama, Higashiuratnachi, Chita Gun, Aichi, Japan, and Stardust xe2x80x98Firebirdxe2x80x99 reliably reproduced true to type. A long testing period followed.\nIn February of 1993, the plants of Stardust xe2x80x98Firebirdxe2x80x99 continued to display the desirable characteristics selected for. These plants were then examined for their susceptibility or immunity from known diseases, namely, Dendrobium virus, Dendrobium Thabdo virus and Tomato Spotted Wilt virus. Stardust xe2x80x98Firebirdxe2x80x99 displayed immunity to those diseases. All developmental work was done in a glass greenhouse located at Aichi-ken, Japan where minimum temperature was kept at 13xc2x0 C.\nStardust xe2x80x98Firebirdxe2x80x99 is descended from 5 species of the genus Dendrobium. These are Dendrobium unicum, Dendrobium moniliforme, Dendrobium aureum, Dendrobium signatum, and Dendrobium nobile. These species have been extensively used in hybridizing and are found over an extremely large area having a wide range of habitat elevations. The habitat area extends from India in the West to Japan in the north, and Australia and New Zealand in the south. Eastward, the habitat includes most of the Pacific Islands. The species are found from sea level to elevations as high as 12,000 feet. It is impossible to make accurate generalizations about the horticulture of the Dendrobium species but it can be said that its hybrids have cultural requirements which are similar to parents Dendrobium unicum and Dendrobium ukon, neither which is patented, and form a predictable guide. Dendrobium unicum and Dendrobium ukon and Stardust xe2x80x98Firebirdxe2x80x99 have rather large attractive flowers and are generally easy to cultivate.\nDendrobium unicum is a very distinctive species from Thailand and Laos. It has a deep orange color. It differs significantly from Stardust xe2x80x98Firebirdxe2x80x99 in that the petals and sepals are very narrow and strongly reflexed, the lip is larger than petals and sepals, is brownish orange and has a 3 ridged callus along the center. Dendrobium ukon is a hybrid, registered with The Royal Horticultural Society in 1979 by H. Furuse. It differs from Stardust xe2x80x98Firebirdxe2x80x99 in growth habit, in the pseudobulbs being quite large and swollen in the middle, and in flower shape and size. Flowers are usually white or pink, petals larger than sepals, with a broad flat lip. Petals and sepals are also much wider at the apex than at the base."}
{"text": "1. Technical Field\nThe present disclosure relates to a light-emitting device such as an organic electroluminescence element (hereafter referred to as “an organic EL element”).\n2. Description of the Related Art\nOrganic EL elements have been examined for application to various light-emitting devices and have been practically used as light-emitting devices. For example, organic EL elements are used for flat panel displays, backlight units for liquid crystal displays, and light sources for illumination."}
{"text": "1. Field of the Invention\nThis invention relates to electrical energy storage devices and is directed more particularly to an electrically rechargeable, reduction-oxidation (REDOX) type cell.\n2. Description of the Prior Art\nREDOX cells utilizing anode and cathode fluids separated by an ion permeable membrane and employing REDOX couples as the electrochemical active materials are generally well known. The electrochemical capacity of these fluids is a function of the amount of active material in the solution and the oxidation state of the material. If the anode fluid and the cathode fluid contain the same number of equivalents of active material and both have the same degree of charge the REDOX system is said to be in balance.\nFor a REDOX system using a chromous/chromic couple as the anode fluid and the ferrous/ferric couple as the cathode fluid, a balanced system has the same number of equivalents of chromous chloride in the anode fluid as there are equivalents of ferric chloride in the cathode fluid. The anode fluid is an aqueous solution of hydrochloric acid and chromium chloride while the cathode fluid is an aqueous solution of hydrochloric acid and iron chloride.\nDuring the operation of a REDOX system, electrochemical imbalance can gradually build up if reactions other than the desired electrochemical REDOX reactions occur. For example, if a system when new had a capacity of 100 ampere hours and after a certain length of time had a capacity of only 90 ampere hours, it would be clear that an imbalance between the anode and cathode fluids had occurred. This imbalance is defined as a different state of charge in the anode and cathode fluids.\nAn imbalance in a REDOX system can occur for a number of reasons including, but not limited to, the following:"}
{"text": "The present invention relates generally to wire bonded integrated circuit packages and, more particularly, to methods and arrangements for isolation and reinforcement of individual bonding wires in integrated circuit packages.\nIn a wire bonded integrated circuit package, a die is attached and wire bonded to a lead frame (or other substrate) and then molded. During the molding process, the flow of the mold may displace some bonding wires, a phenomenon commonly referred to as “wire sweep”. As a result of wire sweep, some displaced bonding wires may come in contact with other bonding wires and produce short circuits. Wire sweep generally causes a decrease in package yields, since packages with shorts are rejected (“short rejects”). In other situations, wire sweep may push adjacent bonding wires closer together than is desired even though they do not actually touch. This can lead to increased inductive coupling and/or crosstalk interference between adjacent bonding wires, which can decrease the performance of the resulting packaged integrated circuit and lead to further reductions in package yields.\nOne approach to addressing the wire sweep problem has been to insulate the bonding wires before molding. By way of example, U.S. Pat. Nos. 5,455,745 and 5,527,742 describe methods of insulating bonding wires. Although the existing wire insulation techniques work well, in the semiconductor industry, there are continuing efforts to improve packaging techniques and/or reduce costs."}
{"text": "An image detection apparatus that utilizes a flat panel detector is widely used in x-ray digital imaging to generate digital image data. The flat panel detector includes a scintillator for generating fluorescence when exposed to radiation and a pixel array for detecting the fluorescence of the scintillator. The scintillator and the pixel array are disposed on, for example, individual substrates, respectively, and are bonded to each other through an adhesive layer. The radiation is converted into light by the scintillator and the fluorescence of the scintillator is converted into an electrical signal by the pixel array, and, as a result, digital image data is generated."}
{"text": "1. Field of the Invention\nThe present invention generally relates to the field of fiber optical circulators. More particularly, the present invention relates to the field of compact, low cost and high performance optical circulators with at least three ports, which can expanded to more ports.\n2. Description of the Prior Art\nIn general, the fiber optical circulator transfers light from a port to another sequentially. It has at least three ports. The light enters the circulator through a first port and exits through a second port. Another light that enters the second port exits through a third port and so on. Since the optical circulator is an inherently non-reciprocal device, the light never goes to other ports.\nSuch an optical device has a broad range of applications in the field of fiber optic communication systems, such as in bi-directional communication, dense Wavelength Division Multiplexing (WDM) and Dispersion Compensation.\nMany types of optical circulators have been developed and manufactured. However, the high cost and bulky size of the circulator have limited the popular acceptance for system applications."}
{"text": "Portable radios with earphones have become popular for use during many types of activities. For example, exercisers use portable radios with earphones to listen to radio broadcasts while exercising. Sports fans use portable radios at live games to listen to radio broadcasts of the game while viewing the game from the stands. Moreover, others use portable radios with earphones to listen to broadcasts privately and therefore, not disturb others. Portable radios are configured so that large headphones which substantially cover a user's ears are positioned adjacent to the user's ears by supporting them on a substantially rigid head mount. Smaller headphones are anchored inside the user's ear canals.\nIn portable radios with headphones, the antenna for the receiver is the long wires attaching the receiver to the headphones. Long antenna wires are required for quality reception of FM transmissions. Frequently, the long wires act as encumbrance to the user of the portable radio. For example, when the receiver is attached to the user's belt, suspended from the user's shoulder by a shoulder strap, or in the user's hand, the length and position of the antenna wires causes them to become tangled and unmanageable. Frequently, the wires are in the way of the user's hands and can become tangled with other articles carried by the user such as a purse or jacket.\nAnother type of portable radio with headphones is that having the receiver fixably attached to the headphone's substantially rigid head mount. In this arrangement, a rigid antenna is also attached to the head mount and typically projects upwards. While this type of portable radio avoids the problem of loose long antenna wires, it is a bulky and unbalanced arrangement, and therefore not suitable for use during activities involving substantial motion.\nOther portable radio configurations include those which run an antenna through a headband so that it encircles a user's head. While this configuration avoids the problem of loose long antenna wires, the headband is rigid and therefore, uncomfortable. In this configuration, a headband made from a more comfortable non-rigid material such as fabric is unacceptable for many activities, such as exercising, because with the wires inside it, the headband cannot be washed.\nOther portable radio configurations also include those which are attachable to a washable headband by clips. Because the radio is removable from the headband, the headband is washable. However, in these configurations, the earphone wires which double as antenna wires hang from the radio, drape down past the user's shoulders and therefore the wires are in the way of the user's hands and can become tangled with other articles carried by the user."}
{"text": "The invention relates to containers, and more particularly to an insulated container possessing impact cushioning characteristics with an improved ability to maintain a desired limited cold or hot temperature range for an extended period of time.\nThe shipment and handling of various goods and commodities, such as blood, drugs, vaccines, frozen foods, catalysts, chemicals, and the like, requires that the shipping container be well suited for cold or hot storage and that it have the ability to maintain the shipped product in a particular environment over a limited temperature range during the course of shipment and handling. In designing such a shipping container, the weight of the empty container is of significance; for instance, many perishable goods travel by air to distant destinations. The size and configuration of the container is also a design consideration, whether ground or air transportation is employed. The materials from which the container is constructed constitute a further factor in its design, since it is desirable to form the container from low cost materials which yield an easily assembled container which, where required, can be quickly disassembled for storage and/or re-shipment. It is further desirable to fabricate the container from materials which are relatively resilient to the customary impact shocks and blows encountered in shipment, and which provides satisfactory protection from damage incurred by the container through rough handling. A shipping container for cold or hot storage must therefore maximize the protection and preservation ability of the container, and minimize the cost and difficulty of construction of the container.\nIt is well known in the prior art to provide a shipping container with foam insulation as both the means to preserve a cold temperature storage environment within the container as well as protect the enclosed product from damage. In Cline et al, U.S. Pat. No. 3,4l6,692, a cold storage shipping container is disclosed which is semi-rigid and self supporting. The container has walls composed of a flexible and resilient foamed insulating material of an interconnected open-cell type which is enclosed within a fluid impervious casing, possessing a similarily composed cover. Since the walls and bottom of the container are an integral unit, the container cannot be knocked-down for storage or shipment. The container is typical of those types of shipping containers having walls composed of an insulating material encased between outer and inner walls. Such a container, though having flexible walls, is relatively expensive, since the encased insulation assembly is fairly elaborate. Further, this type of container must be purchased in assembled form, and cannot be compactly shipped or stored by the user for later assembly, as needed.\nLikewise, it is known in the art to use an open-celled insulating foam material for insulating and protective purposes in containers. Such foam material as used in Cline, U.S. Pat. No. 3,416,692, is of an interconnected open-cell type, such as polyurethane, where the interstices between the foam cells reduce the circulation of gases from the colder interior compartment through to the outer warmer side of the container. However, the ability of such open-celled foam to hold a reduced temperature within a shipping container is far inferior to closed-cell foam such as polystyrene, which material has a relatively low susceptability to internal convective circulation or direct transmission of gases.\nErnst, et al, U.S. Pat. No. 3,890,762, shows a produce shipping container comprised of a common corrugated paperboard box and six closed-cell polystyrene foam panels, the six panels respectively residing in facial engagement with the sides, top, and bottom of the box. The panels are of a rigid foam, however, and thus provide a limited amount of protection to the contained product from damage through shipping and handling. Rigid panels of this type are also subject to chipping and breaking, which effectively degrades the container's ability to maintain a reduced interior temperature, and limits the commercial life of the container itself for re-use where desired.\nThe temperature preservation ability of the Ernst container further suffers from the employment of a type of closed-cell foam which will conduct the gases in the interior of the container to the ambient atmosphere or vice versa. The absence of an effective barrier to the flow of vapors between the interior and exterior of the container significantly reduces the shipping container's ability to maintain the requisite storage environment. In point of fact, the foam material specified in Ernst was selected for its ability to permit rapid temperature modification of the contained produce through cooling by an external source.\nThe prior art thus fails to provide an inexpensive and easily assembled light-weight, reusable shipping container having the improved ability to maintain a desired limited cold or hot temperature range over a long period of time while providing adequate cushioning and protection for the contained product."}
{"text": "A projection system typically includes a light source, illumination optics, one or more imagers, projection optics, and a projection screen. The illumination optics collect light from the light source and direct that light in a predetermined manner to the imager(s). The imager(s) are typically polarization-rotating, image-forming devices, such as liquid crystal display imagers, which operate by rotating the polarization of the light to produce an image corresponding to digital video signals. The projection optics then magnify the image and project it onto the projection screen.\nImagers used in projection systems typically rely on polarizers to separate light into a pair of orthogonal polarization states (e.g., s-polarization and p-polarization). The separation generally involves transmitting light having a desired polarization state and absorbing the majority of the light having an undesired polarization state. Such polarizers, however, absorb a substantial amount of light energy, which is converted to heat. During extended periods of use, the generated heat may build up and damage the polarizers, thereby diminishing the polarization properties. This accordingly reduces the quality of the projected images over time. Typically, high-speed air flows are provided over the polarizers to maintain the polarizer temperatures below acceptable limits. However, the required volumetric flow rates may be high, and the size and number of fans needed in the projection system can consume a significant amount of electrical power and space in the projection system, in addition to generating distracting noise."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a synchronization method of communication, in particular, to a synchronization method adapted to be employed in the third generation mobile communication system for decreasing the call-dropping probability thereof.\n2. Description of Related Art\nWith the arrival of the information era, demands for mobile audio, data communication, and kinds of mobile service are continuously increasing. Limited communication channels and insufficient transmission speed deter the concurrent communication system to be further developed. Therefore, the third generation mobile (3G) mobile communication, capable of providing higher frequency band utilizing efficiency and faster transmission service, has been developed accordingly.\nFIG. 1 is a schematic diagram for illustrating a 3G mobile communication network. Referring to FIG. 1, a user equipment (UE) 10, e.g., a cellular phone or one of other handheld communication devices, communicates to a universal mobile telecommunication system radio access network (UTRAN) 11 via radio signals. The UTRAN 11 includes several radio network subsystems (RNS) 121, 122. Each of the RNSs is controlled by a radio network controller (RNC) 131, 132. Each RNC is connected with a plurality of node bases (Node B) 141-143. For example, as shown in FIG. 1, the RNS 121 includes RNCs 131 and Node Bs 141, 142, and the RNS 122 includes RNCs 132 and a Node B 143. The UTRAN 11 communicates with other telecommunication network, e.g., internet, via a core network (CN) 15 to provide 3G audio and/or data transmission service.\nIf moved or the wireless environment is changed, the UTRAN may shift UE among different RNSs, e.g., changing from connection to the RNS 121 to connect to the RNS 122. Such a procedure is called serving radio network subsystem (SRNS) relocation. In such a procedure, the UE 10 and the UTRAN 11 have to communicate to each other by synchronously using identical setting configuration. For example, the UE 10 and the UTRAN 11 must synchronously use data cipher corresponding to a new setting configuration to encrypt or decrypt the data.\nFIG. 2 is a schematic diagram for illustrating the 3G mobile communication synchronization between a UE and a UTRAN. Referring to FIGS. 1 and 2, in an SRNS relocation procedure, the UTRAN transmits a radio resource control (RRC) message containing ciphering configuration information. After receiving the RRC re-configuration message containing the ciphering configuration information, the UE changes to the corresponding new ciphering configuration and transmits RRC response message to the UTRAN. For example, the UTRAN transmits a UTRAN mobility information RRC message containing downlink counter synchronization information to the UE. After receiving the aforementioned RRC message, the UE changes to the new configuration and feeds back a UTRAN mobility information confirm RRC message to respond the reception of the RRC message.\nBecause of the standard of 3G mobile communication, the UE is setup to measure e.g. wireless channel, according to which the RNC performs and controls channel source management. The obtained results also depend on the adjustment of the wireless channel quality factors for maintaining quality of different services. Therefore, the UE transmits a measurement report at predetermined times to the UTRAN according to the measurement event or the standard. However, the 3G mobile communication standard does not clearly define when the UTRAN should use its corresponding new ciphering configuration. As such, the UTRAN cannot correctly use the same data cipher of the UE to encode or decode the data, which is illustrated with reference to FIG. 4 as follows.\nFIG. 3 is a schematic diagram for illustrating incorrectly encrypting/decrypting due to a synchronization failure between the UE and the UTRAN of 3G mobile communication. At the beginning, the UTRAN transmits UTRAN mobility information RRC message containing downlink counter synchronization information and ciphering mode info containing new ciphering configuration to the UE. However, after receiving the aforementioned RRC message, the UE most probably sends measurement reports to the UTRAN at predetermined times triggered by measurement event or following the standard, before changing to the new ciphering configuration and feeds back a UTRAN mobility information confirm. In such a way, the UE will use the old configuration to transmit the measurement report. Unfortunately, if the UTRAN has already used the new ciphering configuration to receive the measurement report, the UE and the UTRAN is unsynchronized, thus causing incorrectly data encrypting/decrypting of the two parties. Affections of the incorrectly data encrypting/decrypting are illustrated below with reference to FIG. 4.\nFIG. 4 is a schematic diagram for illustrating incorrectly data encrypting/decrypting of the UE and the UTRAN. Referring to FIG. 4, it is assumed that the measurement report transmitted by the UE is composed of protocol data units (PDUs) 401-402, while the UTRAN mobility information confirm is composed of PDUs 403-404 of the RLC entity. If the data are asynchronously encrypted/decrypted, the UTRAN would decrypt the received PUDs 405-408 into incorrect service data unit (SDU), e.g., SDU 411 and SDU 412 which would not form RRC messages. The measurement report includes measurement values of the wireless connection, and the UTRAN can select exchanging cells according to the measurement values. Most importance of all, the new configurations of data cipher between UE and UTRAN are not synchronous any more and the old configuration cannot be changed back. Therefore, the call-dropping probability is increased."}
{"text": "The present invention relates to a novel method for treating organic wastes, and more specifically relates to a method for producing a feed by using organic wastes and oils as the raw materials to form a koji-fermented feed.\nAdditionally, the present invention relates to a method for treating highly concentrated organic wastes with a BOD of at least 10,000 ppm such as sake lees and food scraps, and further relates to a crude fiber feed comprising a ligneous matrix obtained by treatment with the above-described treatment method.\nIn recent years, the problem of how to treat waste from the food processing industry or highly concentrated organic waste such as food scraps has become one which has bearings on the public health. This is because they contain large amounts of water and are perishable, generating malodorous fumes in the process. While this type of waste is usually simply discarded, the burial of these substances results in their decay in the soil, producing large amounts of nitrate-nitrogen which can contaminate ground water. Additionally, if incinerated, they can generate dioxins and NOX, which can pollute the atmosphere. Moreover, the amount of such organic waste is increasing annually at an alarming pace.\nTherefore, many methods of fermenting organic waste to make fertilizer or methods for treating waste for decomposition in order to enable larger amounts of waste to be collected and recycled have been proposed in recent years, but with conventional methods, the treatments can take months, and their utility has been limited due to the risk of saline pollution when used in large amounts because of the salt content in the resulting fertilizers. Furthermore, such waste often has a high water content, making storage and transport difficult. As a result, this type of waste has not been able to be used effectively.\nOn the other hand, there are methods for making such raw garbage into fertilizer using microbes, but such microbes are not very compatible with garbage having a high oil content, so that there are limits to the range of raw waste that can be treated in this manner.\nThe present invention has been made in consideration of the above-described circumstances, and has a first object of offering a novel means for making effective use of the above-described organic wastes.\nThe present invention has the further object of offering a method for producing a feed by efficiently processing the above-described organic wastes.\nThe present invention also has the object of offering a feed produced by the above method.\nThe second object of the present invention is to offer a method for treating organic wastes such as sake lees and food scraps, especially highly concentrated organic wastes with a BOD of at least 10,000 ppm, with extreme efficiency in a short period of time.\nAnother object of the present invention is to offer a method for treating an organic waste capable of converting a large amount of waste into a readily disposable form, and use of a substances treated in this way.\nAs a result of diligent research towards achieving these objects, the present applicant discovered that if oils are added when treating organic waste with certain molds, the molds can transform the oils into fuel and efficiently convert them into the heat of fermentation. That is, they found that by employing oils as energy sources for the molds to dramatically improve the decomposition and drying speed of the organic matter due to the mold, it is possible to efficiently treat organic wastes. Additionally, more surprisingly, they discovered that when ammonia gas and acetic acid are generated, the generation thereof can be prevented by mixing in some oils, that is, the addition of oils can inhibit the proliferation of putrefying bacteria such as acetic bacteria.\nThe present invention offers a method for treating organic wastes, comprising steps of mixing said organic wastes with a mold to obtain a mixture, and allowing said mixture to ferment and dry, wherein the oil content of said mixture is at least 2 wt %, preferably at least 5 wt %, more preferably at least 10 wt % with respect to the water content of said mixture. In this way, the heat of fermentation of the mold can be used to efficiently dry the wastes.\nAccording to a first object of the present invention, a feed is produced by adjusting the oil content of a mixture of organic wastes and koji mold or koji to at least 3 wt %, preferably at least 5 wt %, more preferably at least 10 wt %, and allowing said mixture to ferment and dry. Preferably, the feed is produced by mixing together organic wastes, oils, koji mold or koji, and allowing the koji mold to ferment and dry the waste.\nBy adding oils as nutrients in this way, the processing speed of organic wastes, in particular highly concentrated organic wastes can be highly improved.\nFurthermore, according to this type of method, raw garbage, livestock excrement and the like which has conventionally been discarded can be effectively used.\nFurthermore, since the method of the present invention makes use of koji mold instead of treatments using high temperature microbes, so that while a portion of the organic matter is decomposed into carbon dioxide and water, the most part is synthesized into microbial proteins, while large amounts of digestive enzymes are secreted for digestive consumption. As a result, a feed of high nutritional value which is suitable for feeding to livestock and the like is obtained. That is, while the treatment of organic waste by microbes has conventionally been for the purpose of breaking down wastes, the method according to the present invention enables the generation of active ingredients by the koji mold, thereby changing the organic waste into feed which can then be put to effective use.\nIn the present invention, oils are used in order to produce a feed fermented by koji mold from organic wastes. Conventionally, the fermentation and heat generation due to koji mold has been mainly used for decomposition of starches, but in the present invention, the oils are used to assist largely in the fermentation and heat generation. Since the koji mold can consume the oils and efficiently convert them to the heat of fermentation, there is no particular need to add nutrients for growth of the koji mold, and the oils act as an energy source to dramatically increase the speed of decomposition and drying of the wastes. Therefore, there is no need to specially add starchy materials even to wastes such as raw garbage with a low starch equivalent. Furthermore, since there are few microbes that can consume oils in nature, the koji mold, which is relatively susceptible to contamination with other microbes, can be grown with dominance. The adjustment of the oil content should preferably be done by intermixing oils, the amount of oils normally added being preferably at least 1.5 wt %, more preferably at least 5 wt %, and most preferably at least 10 wt % of the water content of the mixed organic waste.\nThe oils used in the present invention can be of any type as long as they are able to achieve the object of the present invention, and can be chosen from among animal, vegetable or mineral oils, or combinations thereof. Some examples are those which are readily available, including edible oils such as frying oil, lard, soy oil, waste oil from meat processing or olive oil, or cutting oils, with those having a boiling point of at least 100xc2x0 C. being most preferable. When preparing a feed, edible oils are preferable.\nAdditionally, the mold in the present invention refers to filamentous eukaryotic microbes (molds), particularly Aspergillus, Monascus, Mucor and Rhizopus, of which Aspergillus and/or Monascus are chosen with preference for having amylase and lipase activity, further among which Aspergillus oryzze kawachii and Aspergillus awamori kawachii (available from Kirishima Kogen Beer KK) have a high oil-decomposing ability and can be used effectively. Normally, oils are avoided due to the inability of conventional fertilizer plants to process them. However, by using this type of koji mold, which has a high oil decomposing ability, a waste with a high oil content can be effectively treated. Additionally, they can be used in combination with other molds which can suitably assist in the decomposition process.\nAccording to a second object of the present invention, an organic waste and oils are preferably added to a matrix material to form a mixture, and a mold is grown in the mixture to treat the organic waste by fermentation thereof.\nIn this case, the organic waste should preferably be provided in an amount such that the water content is 20-70 wt %, preferably 20-40 wt % with respect to the matrix material.\nAdditionally, the matrix material used in the treatment method of the present invention can be of any type as long as it is capable of achieving the object of the present invention, but should preferably be a pulverized ligneous material. For example, sawdust, wood shavings, and their mixtures are suitable for use.\nAdditionally, the mold can be allowed to occur naturally in the mixture, or positively added to the mixture. The amount added should be in the range of about 0.01% to 0.1% of the total dry weight of the mixture.\nAccording to the present invention, the treatment process is usually a continuous process such that as the organic waste is being treated, new batches of the organic waste and oils can be added as needed. Here, the matrix material can also be supplemented if necessary.\nIn the processing method according to the present invention, the mixture is preferably adjusted to 50xc2x0 C. or less, preferably about 40-50xc2x0 C. This temperature regulation should preferably be performed by ventilation, and it is especially convenient to blow warm dry air which is available from the exhaust heat from the generator of the ventilator.\nAdditionally, once the koji mold begins to grow and fermentation begins, the heat of fermentation causes the waste to dry so as to result in a feed with a water content below a desired level, and this can be improved by making proper adjustments to the speed and level of drying due to the ventilation of warm air. The final water content of the feed should preferably be about 30 wt % or less, preferably 20 wt % or less, and most preferably about 15 wt % or less.\nFurthermore, in the present invention, a feed obtained by the above-described processing method is offered. This feed has an extremely high nutritional value, and is suitable for feeding to livestock such as cattle, swine and fowl. Furthermore, when using a matrix material, a matrix material obtained by treating organic waste can be used as a crude fibrous feed. Since the koji mold uses oils as an energy source, the oils are consumed and reduced. Moreover, the feed will contain lipases from the koji mold. While oils can cause diarrhea in livestock, this problem can be completely overcome due to the presence of lipases, so that the oils can actually be used to raise the nutritional value of the feed and promote the growth of the livestock. Furthermore, the action of the koji mold also leaves large amounts of other enzymes such as amylase, protease and active oxygen-decomposing enzymes in the feed, which can then largely assist in digestion by livestock to which the feed has been fed, thereby relieving stress. Additionally, it also has the effect of suppressing the foul odor of livestock excrement.\nFurthermore, while livestock fed feeds containing unsaturated fatty acids can have reduced meat quality due to yellowing or the like, treatment with koji mold reduces the amount of unsaturated fatty acids, thus giving the livestock fed the feed of the present invention good meat quality."}
{"text": "This invention relates to trans(.+-.) or trans(-)-octahydropyrimido[4,5-g]quinolines of the formulas ##STR1## wherein R is H, CN and C.sub.1 -C.sub.3 alkyl or allyl; R.sup.1 is H, OH, SH, O--C.sub.1 -C.sub.3 alkyl, S--C.sub.1 -C.sub.3 alkyl or NH.sub.2 ; R.sup.2 is H, OH or NH.sub.2 with the proviso that R.sup.1 can only be NH.sub.2 when R.sup.2 is OH or NH.sub.2 ; and acid addition salts thereof made from non-toxic acids. Formula I is named as a trans-(.+-.)-6-substituted-2,4-permissibly-substituted-5,5a,6,7,8,9,9a,10- octahydroyrimido[4,5-g]quinoline and formula II is the corresponding trans-(-)-enantiomer.\nPreferred subgroups of compounds represented by I or II above include:\n(a) compounds in which R is C.sub.1 -C.sub.3 alkyl;\n(b) compounds in which R is n-propyl;\n(c) compounds in which R.sup.2 is H; and\n(d) compounds in which R.sup.1 is SH or S--C.sub.1 -C.sub.3 alkyl.\nThe compounds represented by formulas I and II absorb ultraviolet light with absorption peaks in the range 265-350 mu, (which range encompasses the UVA--280-300 mu range--) and are therefore potentially useful as sunscreen agents. The compounds are also useful as intermediates, as will be more fully developed hereinafter.\nCompounds according to I or II in which R.sup.1 is SH or S--C.sub.1 -C.sub.3 alkyl and R.sup.2 is H are prepared by reacting thiourea or an S--C.sub.1 -C.sub.3 alkyl isothiourea with a compound of the formula III or IV ##STR2##\nCompounds according to I and II wherein R.sup.2 is H and R.sup.1 is OH are prepared by acidic hydrolysis of a trans-(.+-.) or trans-(-)-2--C.sub.1 -C.sub.3 alkylthio-6--C.sub.1 -C.sub.3 alkyl or allyl-5,5a,6,7,8,9,9a,10-octahydropyrimido[4,5-g]quinoline. The product of this reaction, a 2-OH derivative, can be readily alkylated by standard etherification reactions; e.g., dimethylsulfate and base, to yield the 2--C.sub.1 -C.sub.3 alkoxy derivatives; e.g., a 2-methoxy derivative, according to I or II.\nCompounds according to I or II in which R.sup.1 is NH.sub.2 and R.sup.2 is OH are prepared by reacting a 1-alkyl or allyl-6-oxo-7-alkoxycarbonyldecahydroquinoline (furnished by the procedure of Schaus, Huser and Booher, Ser. No. 535,519, filed this even day, whereby a trans-.+-. or trans---6-oxo-1--C.sub.1 -C.sub.3 alkyl or allyldecahydroquinoline is reacted with a dialkyl carbonate in the presence of sodium hydride in an inert mutual solvent, with a guanidine salt. After work-up of the reaction mixture, the 2-amino-4-hydroxy derivative thus formed is isolated by standard means.\nCompounds according to I or II in which both R.sup.1 and R.sup.2 are NH.sub.2 are prepared by cyclizing the usual trans-(.+-.) or trans-(-)-1-substituted-6-oxodecahydroquinoline with cyanoguanidine in a high boiling mutual inert solvent such as CARBITOL--see Modest et al., J. Org. Chem., 30, 1837 (1965). When R.sup.1 and R.sup.2 are H, a different cyclization reaction is employed to join the pyrimidine ring to a 1--C.sub.1 -C.sub.3 alkyl or allyl decahydroquinoline nucleus. The procedure of Boger et al. J. Org. Chem., 47, 2673 (1982) was adapted to the specific synthetic goal. This procedure involves the reaction of a cyclic ketone with pyrrolidine to form an enamine (Formula V or VI below), reaction of which with 1,3,5-triazine yields an unsubstituted pyrimidine ring fused to the decahydroquinoline. ##STR3## wherein R is C.sub.1 -C.sub.3 alkyl or allyl.\nAn illustrative reaction from Boger is that of cyclohexanone with pyrrolidine to form a 1-pyrrolidino-1-cyclohexene which, on treatment with 1,3,5-triazine yields tetrahydroquinazoline. In our reaction, the trans-(.+-.) or trans-(-)-1--C.sub.1 -C.sub.3 alkyl or allyl-6-oxodecahydroquinoline is the starting ketone which forms the enamine with pyrrolidine.\nWhile compounds according to I and II in which R is C.sub.1-C.sub.3 alkyl or allyl are the primary products of this invention, those compounds in which R is H or CN are not only UV absorbers but are also intermediates for the preparation of compounds carrying a different alkyl group (from that of the 6-oxodecahydroquinoline starting material) or an allyl group. For example, if it is desired to prepare a compound according to I or II in which R is allyl, it is possible (depending on the nature of the R.sup.1 and R.sup.2 substituents) to replace the R alkyl group with CN (using CNBr), and then to remove the CN group by hydrolysis to yield compounds wherein R is H. This secondary amine derivative can then be allylated as with allyl chloride or can be reductively alkylated with an aldehyde or can be alkylated with a C.sub.1 -C.sub.3 alkyl halide to yield a compound according to I or II having an allyl group or a different alkyl at N-6. Incidentally, it is not necessary to replace an N-methyl with a different alkyl group in order for the above process to be useful. For example, it is possible to replace a CH.sub.3 group with a .sup.13 CH.sub. 3 containing an isotopically-tagged carbon. Preferably, however, because of the reactivity of SH and OH groups at C-2 in I or II to CNBr, the interchange of groups on the quinoline nitrogen should take place prior to ring closure; i.e., at the 6-oxodecahydroquinoline stage. In this procedure, a trans-(.+-.) or trans-(-)-1-methyl-6-oxodecahydroquinoline, for example, is transformed to the 1-cyano compound--see Bach et al., J. Med. Chem., 23, 481 (1980) or U.S. Pat. No. 4,198,415. Acidic hydrolysis or hydrogenolysis of this compound by the methods of Titus and Bach, Ser. No. 535,522, filed this even day, yields trans-(.+-.) or trans-(-)-6-oxodecahydroquinoline. This secondary amine can then be allylated or realkylated possibly with a tagged alkyl, using allyl or alkyl halides to prepare the desired intermediate trans-(.+-.) or trans-(-)-1-allyl or alkyl-6-oxodecahydroquinoline. The above procedures are clearly extremely useful in preparing the difficult-to-obtain N-allyl derivatives because the best procedures for obtaining the 2-alkyl-6-oxodecahydroquinolines involves multiple reduction or hydrogenation procedures, and an allyl group would not survive all of these synthetic steps."}
{"text": "In general, there are two kinds of standards to ensure the performance of an audio transducer and associated holding device (e.g., telephone handsets, telecommunication system handsets, etc.) on the audio receiving path. One such standard is based on acoustic performance, the conventional manner of using the handset by acoustic coupling to normal human ears. The other standard is the electromagnetic performance for a human who is hearing impaired or otherwise disabled, and that needs to use a hearing aid. Passing the compliance for both acoustic performance and electromagnetic performance simultaneously is very challenging. This is especially true for passing the acoustic compliance standards with an artificial ear (e.g., a “Type 3.3” standard), which is the closest representation to a human ear currently used that can reflect the way people use a listening device (e.g., handset, etc.)."}
{"text": "Semiconductor components of this type, be they thyristors or diodes, for example, are generally known. The semiconductor body of a thyristor generally possesses four zones of alternate opposite conductivity type, and a pn-junction in each case lies between these zones. The pn-junction located between the outer and the inner anode-side zone is referred to as a blocking-pn-junction, since the greatest part of the blocking voltage occurs across it. The semiconductor body of a power diode basically possesses two zones of opposite conductivity type. The semiconductor bodies of such semiconductor components generally possess a specific content of recombination centers, which decisively influence their electric properties, for example the turn-off time. The aforementioned recombination centers are generally heavy metal atoms, such as gold, platinum or manganese.\nThe processes which are employed in the production of the above-mentioned semiconductor components involve the gettering of a specific part of the introduced recombination centers into regions close to the surface. The gettering layers are, for example, highly doped zones, in particular zones doped with boron or phosphorus, or also zones close to the surface, whose crystal lattice structure is heavily disturbed by mechanical processing. As a result of the gettering, there is a reduction in the concentration of the recombination centers beneath the gettering layer close to the surface. If uniformly gettered layers are present both at the anode side and at the cathode side, the curve of the concentration of the recombination centers over the thickness of the semiconductor body has a U-shaped, symmetrical profile. If, on the other hand, a weaker gettering effect exists at the anode side than at the cathode side, for example, due to less disturbance of the surface, the concentration of the recombination centers is higher at the anode side than at the cathode side.\nThe absolute degree of the concentration of the recombination centers can be set by the quantity of the diffused-in recombination centers, for example, of gold atoms. The symmetry of the profile of the concentration of the recombination centers cannot in practice be effectively influenced either by controlling the quantity of the recombination centers or by controlling the diffusion parameters during the diffusing-in of the recombination centers."}
{"text": "The present invention relates to a system for surveying rails, in particular running rails for cranes, shelf handling units, running wheel blocks.\nA measuring arrangement for monitoring railway tracks is known from DD 212 931 and includes a transmitter unit arranged on the track and having a laser alignment unit. The laser beam is so aligned hereby as to serve as measuring axis, fixed in space, for the track rail. The measuring arrangement includes a measuring vehicle having a traveling mechanism which has two horizontal running rollers as well as two side guide rollers on the right side and two side guide rollers on the left side. The position of the measuring vehicle in relation to the measuring axis is determined by providing a photosensor which confronts the laser beam and is configured as four-quadrant photodiode. Upon incident laser beam, the photodiode generates an electric signal which is delivered to an electronic evaluation device which controls a vertical tracing device as well as a horizontal rotary mechanism for shifting and turning the photodiode. The control is so implemented that each of the four-quadrants of the photodiode occupies in relation to the laser beam the same position. In view of the high sensibility of the photodiode, the positional determination is limited in such a measuring arrangement only by the accuracy of the mechanical adjustment.\nThe drawback of this conventional measuring arrangement resides in the mechanical tracing of the photodiode, which is determinative for the accuracy with respect to positional determination of the rail, i.e. in particular the positional change of the running surface of the rail. Moreover, the long measuring times as a result of the mechanical readjustment of the photodiode in the measuring plane (X-direction and Z-direction) are disadvantageous.\nIt is thus an object of the present invention to provide an improved system for surveying rails, obviating the afore-stated drawbacks.\nIn particular, it is an object of the present invention to provide an improved system for surveying rails, which is capable of approaching individual positions at great accuracy in order to determine the positional change of the running surface of the rail relative to the laser beam, without requiring a shift of the photosensor.\nThese objects, and others which will become apparent hereinafter, are attained in accordance with the present invention by providing a transmitter unit arranged on the rail and emitting a laser beam which is aligned longitudinally in the direction of a rail and is incident on a measuring surface; a detector unit designed for mobility along the rail and having a photodetector which confronts the laser beam and has a matrix in the form of a plurality of neighboring optical sensor elements for generating electric output signals in response to the incident laser beam on the measuring surface; an electronic evaluation unit receiving the output signals from the sensor elements for determining the point of impact and thereby the position of the laser beam on the measuring surface; and a proximity sensor for determining a change in distance between the transmitter unit and the detector unit.\nIn accordance with the present invention, the photodetector may include a matrix of rectangular configuration with a plurality of immediately neighboring pixels (optical sensor elements) for generating the electric output signals which are delivered to the electronic evaluation unit by which the point of impact of the spatially, positionally stabilized laser beam in the measuring surface is determined solely through electronic means by evaluating the output signals of the pixels. The photodetector may be an opto-electronic camera which has arranged in front thereof a transparent diffusion screen as measuring surface which is optically projected onto the matrix of the camera. In this fashion, it is possible to determine for each spot in relation to the length of the rail the positional change of the running surface with respect to the laser beam which is fixed in space. The configuration of the photodetector in the form of a rectangular matrix permits a two-dimensional measurement of the intensity distribution (Gaussian distribution) of the laser beam, thereby realizing a precise determination of the point of impact of the principal beam direction or of the laser beam axis upon the preset measuring surface. Thus, there is no longer any need to shift the measuring surface. The adjustment of the zero point (point of reference) is necessary only once, thereby shortening the duration of the measurements and, moreover, enhancing the measuring accuracy.\nSuitably, the photodetector may be configured as CCD camera. This results in a rail surveying system which is very easy to compile.\nA more precise determination of the point of impact of the laser beam upon the measuring surface can be realized by a so-called subpixel resolution, through determination of the point of impact of the laser beam upon the diffusion screen by conforming the measured subpixel position of the beam image of the diffusion screen upon the matrix. This is attained by initially evaluating each image line and each image column separately in a square zone about the beam image, and subsequently averaging the outcomes. The result provides the precise position (subpixel position) of the beam within the camera image. The difference between the inputted starting position (reference point) is thus a measure for the deviation of the mobile detector unit from the ideal line. Suitably, the conformation is implemented by a computer program.\nAccording to another feature of the present invention, the electronic evaluation unit includes a microprocessor to implement in a simple manner complex computations and control processes.\nIn order to determine the position of the running surface of the rail in dependence on distance, the proximity sensor may include a friction wheel which rolls on the running surface of the rail. Suitably, the friction wheel is connected with an incremental shaft encoder for determination of the travel distance of the detector unit, thereby realizing a high distance precision. The friction wheel may also be utilized to register a state of motion of the detector unit to thereby ensure a correct execution of the survey.\nSuitably, the detector unit can be remote-controlled so that rails which are difficult to access can be surveyed even under adverse conditions of visibility and unfavorable conditions of the surrounding (toxic gases, vapors etc.).\nThe positional stabilization of the laser can be realized by an electronic level which accomplishes a high positional stability in a fairly simple manner. Suitably, the positional stability amounts to +/xe2x88x921 mm deviation over 100 m rail length.\nThe precise approach of the measuring positions can be realized for a running rail cross section composed of flange parts and web parts, by providing at least two vertically rotatably supported guide rollers which are arranged in spaced-apart relation longitudinally in the direction of the rail, and which are forced against both side surfaces of the rail for guiding the detector unit, and by providing at least two running rollers which are arranged in spaced-apart relation longitudinally in the direction of the rail and rolling on the running surface of the rail and which are horizontally supported for rotation about pivot axes aligned transversely to the longitudinal rail direction, with at least one of the running rollers being driven.\nSuitably, the vertical height of each of the guide rollers can be individually adjusted so that a survey of heavily worn rails can be implemented.\nThe traveling behavior of the detector unit on the rail can be enhanced, when, as viewed in longitudinal rail direction, the leading running roller is arranged behind the leading guide rollers. Also, the stability of travel can be improved, when the distance between the axis of the running roller and the axis of the immediately neighboring guide roller is one time to three times the guide roller diameter.\nIn order to ensure a precise movement of the detector unit, the guide rollers on one side of the rail have a fixed pivot axis and the guide rollers on the other side of the rail are forced, e.g. by a spring force, upon the side surface of the rail. The guide rollers may be supported in ball bearing, in particular, self-aligning ball bearings, to better conform the guide rollers to inclined rails.\nIn order to be able to travel over expansion joints and cracked rails, the guide rollers are arranged in pairs on a common transverse carrier which extends longitudinally in the direction of the rail and is articulated to the detector unit. At the same time, horizontal forces, which act during acceleration and deceleration, are thereby spread over more contact points so that a more stable guiding behavior is realized.\nAccording to another feature of the present invention, the running rollers have a running surface which is formed in mid-section with a radial groove to realize a better stability with respect to lateral inclination as a consequence of the attained two-point contact upon the rail surface.\nA safeguard of the detector unit can be implemented by providing a safety system. Defective rails or rail sections, which can lead to a crash of the detector unit, are registered in particular by providing the safety system with at least one proximity sensor for determining the distance of a reference point of the detector unit relative to the rail.\nA survey of at least two rails extending in one plane can be realized by splitting the laser beam into two horizontal partial laser beams extending perpendicular to one another, for determination of the height difference of both rails with respect to a horizontal straight line extending in a vertical plane. This enables a precise measurement of the parallelism of both rails relative to one another in space in a simple manner. Suitably, the partial laser beam, which extends transversely to both rails, is positioned in the vertical plane, i.e. the vertical plane is unambiguously defined by this partial laser beam. Advantageously, one of the options for defining the reference line permits the use of the partial laser beam, which extends transversely to both rails, for formation of the horizontal straight line which is oriented in a vertical plane."}
{"text": "In the art of arc welding, particularly in the area using inert gas as a shield, (commonly known in the art as metal inert gas or xe2x80x9cMIGxe2x80x9d systems), generated heat is extensive, being in excess of 10,000 to 12,000 degrees Fahrenheit, at the arc point. It is not uncommon for heat to be transferred to metal parts of the welding apparatus in excess of 1,000 degrees Fahrenheit. There have been multiple ongoing efforts to remove generated heat from these type of welding units and similar units. Excessive heat makes workers uncomfortable, but it also degrades the working apparatus much more quickly during high temperature heat exposure. Replacement of the arc welding apparatus parts that degrade is an expensive and time-consuming process.\nWhen robotic welding units are used, the loss of time required for maintenance becomes significant. For this reason and others, efforts to remove the built-up heat generated during the welding process has been ongoing, but thus far no system has realized a truly efficient process.\nSome of the earlier attempts at cooling the welding apparatus is exemplified in U.S. Pat. No. 4,864,099 (Cusick), entitled xe2x80x9cWater cooled semi-automatic welding gunxe2x80x9d. This invention utilized a liquid cooled welding gun assembly. Power, gas and cooling liquid hoses were contained within the cable assembly and intercommunicated the connector block of the welding gun assembly with the rear connector block. This welding gun assembly included a liquid cooled nozzle assembly, with internal cooling passages. Although liquids showed a great propensity for cooling, the amount of water to be used in removing the heat because burdensome, and further improvements were needed.\nReferring now also to U.S. Pat. No. 5,670,073 (Kiilunen), entitled xe2x80x9cWater cooled gas metal arc welding gunxe2x80x9d, a water cooled gas metal arc welding (MIG) gun was disclosed in which a water cooling system for the gun barrel introduced coolant water at the bottom of the barrel to create a coolant water flow through the gun. One of the goals and benefits of this invention was that it did not allow the coolant water to stagnate and heat up. Again, the amount of water and burden of maintenance made this particular welding apparatus problematic for heavy use. Another dangerous use of water in the TIG and MIG welding systems was the significant levels of electricity used in the welding process, which made a water leak extremely hazardous, due to electrical shock.\nWater generally presents other problems to the welding processes, beyond increasing the danger of electrical shock. Water intrusion into the welding process often degrades the welding quality of the finished product. There have been numerous patents directed to avoiding the use of water exposure to the welding process. For example, in U.S. Pat. No. 6,153,847 (Nakatani et al.), entitled xe2x80x9cWelding member and welding methodxe2x80x9d a water-repellent film was applied directly onto the surface of a welding member between two mating welding bevels, a welding bevel of a base member to be welded, and a coated welding rod formed by coating a metallic core wire with a flux. The water-repellent film was formed of a hydrophobic material having a contact angle of 90.degree. or above and providing minute concavities and convexities. Thus, in a highly humid environment, during rainfall or even in water, welding could be effected without weld defects such as blow holes, cracks and the like.\nIn a effort to avoid the problems associated with water cooled systems, further improvements were made. Referring now to U.S. Pat. No. 5,248,868 (Cusick), entitled xe2x80x9cMelding gunxe2x80x9d (corrected title is xe2x80x9cWelding gunxe2x80x9d) Cooling in this instance was accomplished through a plurality of open recesses or channels whereby pressurized air was conveyed. This welding gun assembly also included a pressurized air cooled nozzle assembly. In this invention, pressurized and refrigerated air cooled non-insulated conductor wires along with cooling a rear connector assembly, a cable assembly, a front connector assembly or handle assembly, a conductor tube assembly, and a nozzle assembly.\nFurther improvements were made in which air was used as a heat exchanger. Referring now to U.S. Pat. No. 6,005,221 (Cusick), entitled xe2x80x9cPressurized air cooled tungsten inert gas welding apparatusxe2x80x9d. This invention was directed toward an improved tungsten insert gas (TIG) arc welding process, having a means to remove heat energy using pressurized air, and a method for removing the heat away from the torch head area. A source of pressurized air was delivered though a hose to the torch head, and directed through heat conductive air tube coils within the torch head. As the air moved through the coils, it acquired some of the heat energy within the welding head. The heated air was then vented into an outer jacket, which contains the electrically conductive cables, and an airspace surrounding the cables. As the heated air moved toward a rear connector block, it contacted the length of the cables, and was further able to transfer some of the heat to the cables. This allowed some of the heat energy contained within the pressurized air to be transferred through physical contact to the metal cable, so that as the air moved farther along the length of the metal cable, it transferred more of its heat energy.\nUse of heat exchanging/removal substances directly on surfaces being welded was shown in U.S. Pat. No. 6,168,843 (Hallo), entitled xe2x80x9cMethod for dissipating heatxe2x80x9d. In this invention, a method was shown for increasing the dissipation of heat from one portion of a surface when heat was applied to another portion of the same surface. A composition of a gel or paste with high water content and a thickener of a mineral clay in a colloidal suspension was used. This system had problems associated with the need to make all of the various applications of the gel or paste, and due to the time constraint requirements, was virtually unworkable in a robotic setting.\nThis invention and method process combines the benefits of heat extraction, using not only water, but also with the use of air supplying a volume area usable as a means to transport water that has undergone conversion to steam away from the welding head. This invention avoids problems of the risk of electrical shock, while still using water, or any liquid capable of undergoing a phase change from liquid to gas when absorbing a significant amount of heat energy per unit of liquid during the phase change, as compared to heat absorption per unit without any phase change, to absorb vast quantities of heat.\nWater is a cheap and readily available heat transfer means, which is undoubtedly why it was originally used in the inventions cited above. The problems with water usage, however, drove the inventive ideas of people using these products to turn to air as a means to remove heat. It is clear that a greater quantity of air is required to achieve the same heat removal properties that a much smaller quantity of water is capable of doing, and thus volumes of air needed to be pumped through the system. While remaining a better alternative than liquid water, it still could only provide reasonable cooling.\nWhat was not considered in any of the above inventions, was the incredible heat absorption capacity of water during an evaporation process. The absorption of heat by a substance always produces some change. In addition to thermal expansion, its temperature may rise or it may undergo a phase change in which its physical characteristics are altered without a change in temperature. This is the difference between xe2x80x9cspecific heatxe2x80x9d, which considers the number of calories or heat energy units required to raise the temperature of one gram of a substance one degree Celsius. The amount of heat transferred during changes of phase is called xe2x80x9clatent heatxe2x80x9d. The heat of vaporization for water is extremely large because of the strong attraction between water molecules and the consequently large amounts of energy that must be expended to separate then into the vapor state. For water, the latent heat of vaporization to convert water to steam, where both remain at 100 degrees Celsius is 540 cal/g.\nSteam, by its very nature of being in an excited state, requires a great deal more volume then it\"\"s a previous liquid state. None of the above inventions provide for the transfer of latent heat energy, and are thus restricted to much lower levels of heat energy transfer, since they must rely solely on xe2x80x9cspecific heatxe2x80x9d transfers.\nThis invention takes advantage of the xe2x80x9clatent heatxe2x80x9d transfer, by providing a small quantity of water which is directed through a common atomizer, and then directed into the welding head unit so it can contact those parts of a welding head that absorb great quantities of heat. The atomized water, being presented in a mist form, is able to fill the airspace provided, contact a heated nozzle, and quickly convert into steam which is then able to be injected directly out of the rear of the arc welding nozzle, away from the user. The amount of water used, and thus the amount of steam generated, is extremely small, especially when used in conjunction with other cooling methods. Typically, the user will not even see any appreciable amount of steam exiting the apparatus.\nUsing only liquid water in the prior art forms shown above, are vastly inefficient as compared with the present invention. For example, a cubic centimeter of water, introduced to absorb heat, where the water is 25 degrees Celsius, is only able to absorb 75 calories of heat, before being expelled for cooling. Then, there is the problem of cooling the water, if it is going to be recycled through the system. Contrast the same cubic centimeter of water with the atomizer steam process, and the ability to absorb heat, has the same cubic centimeter of water absorbing the 75 calories, plus an additional 540 calories that is required to convert the water to steam, for a total of 615 calories. This is over eight times as efficient as the liquid water transfer. The efficiency is further enhanced, since the atomized water, having such a small diameter size, is able to undergo the steam conversion virtually immediately when contacting an extremely hot welding head. This makes quicker use of the latent energy transfer, since only small amounts of atomized water are contacting the hot surfaces at any given time, so that they are raised to 100 degrees Celsius almost immediately, and then undergo the acquisition of latent heat, which is the actual efficient energy transfer. By using an atomizer, the efficiency of the system is virtually always using latent heat exchange.\nUsing atomized water allows a mist to be introduced to the heated portions of the welding apparatus. This does not cause temperature gradient stress on the parts of the apparatus, that would otherwise occur if a quantity of liquid were simply splashed onto the part. The instantaneous boiling would produce a problem with the rapid expansion of steam, as well as cause the parts themselves to be subject to cracking and other degradations.\nThe advantages of this innovative system, being able to rely on latent heat transfer is quickly realized when the welding itself is ceased, and the welding apparatus is able to be touched with unprotected human hands within seconds of the termination of the welding. Interchanging parts is able to be done immediately, with the chief benefit, however, being that the parts do not absorb significant amounts heat that causes them to degrade. Therefore, repairs to a welding apparatus are minimized. This system does not produce vast quantities of steam that put other workers at risk, or introduce undesirable amounts of water vapor against other objects.\nIt is therefore the object of this invention to provide an improved welding apparatus that is able to remove undesired quantities of heat energy from a welding apparatus, with the heat transfer being done virtually immediately through use of latent energy transfer.\nIt is also a further object of this invention to provide an improved method of removing undesired quantities of heat energy from a welding apparatus, with the heat transfer being done virtually immediately through use of latent energy transfer with atomizer water vapor converted to steam.\nIt is also a further object of this invention to provide an apparatus and method suitable for robotic, semi-automatic and automatic welding apparatus, that allows continuous use of a welding apparatus with only minimal repairs necessary to the welding apparatus due to heat degradation."}
{"text": "1. Field of the Invention\nThe present invention is related to methods and systems for telecommunications, and in particular, to methods and systems for processing calls.\n2. Description of the Related Art\nConventional corporate enterprise switching systems, such as PBX systems, offer features to make the calling process more convenient. One such feature, termed a callback or campon feature, is used when a caller calls a telephone station that is busy. The callback process determines when the called telephone station is no longer busy and then calls the caller and the called party. In particular, if the called party's telephone line is busy when called, the caller can, via a keypad command, instruct the enterprise telephone system to automatically call the called party's telephone when the called party's telephone line is idle. The telephone system notifies the caller that the called party's telephone is idle by ringing the caller's telephone in a unique manner. If the caller then answers the telephone, the called party's telephone is automatically called and connected to the caller's telephone.\nThe callback feature relieves the caller from the aggravation and wasted time resulting from repeated attempts to place a call to a busy line, hoping that the line will free up. Unfortunately, the techniques used to accomplish the callback feature within an enterprise switching system, such as a PBX system, are not applicable to other telephone systems, such as cellular phone systems. Further, many conventional systems that do provide callback features outside of an enterprise system are very limited in where they can be located. For example, these conventional systems may need be located at a local exchange of the calling party or between the calling party and the local exchange, and typically also are burdened with having to intercept all the callers calls at the point of origination. Further, these types of conventional systems require that the caller be a callback feature subscriber. Other conventional systems simply repeatedly retry the called number in order to determine when the called line is idle, which is an inefficient use of telephony resources. In addition, in order to perform a callback, many conventional systems are burdened with having to perform a call bridging operation. Hence, conventional telecommunications systems fail to offer the callback feature for many telephony applications."}
{"text": "1) Field of the Invention\nThe present invention relates to a technology for notifying presence of a vehicle to a periphery of the vehicle appropriately.\n2) Description of the Related Art\nIn recent years, in place of automobiles that use a gasoline engine (an internal combustion engine) as a power source, hybrid automobiles that include both a gasoline engine and an electric motor as power sources, fuel-cell-powered automobiles that travel while generating electric power with a fuel cell using a hydrogen gas as a fuel, electric automobiles that use a battery (an electric motor), which is charged by a charger, as a power source, and the like have started to be widely used.\nAn advantage in driving a vehicle with such an electric motor is that a traveling sound is small (a silence characteristic is high). However, a pedestrian or a bicycle rider may judge whether an automobile (a vehicle) is present around the pedestrian or the bicycle rider or assume a distance between the vehicle and the pedestrian or the bicycle rider according to an engine sound emitted from the vehicle.\nSince the pedestrian or the bicycle rider makes his/her judgment based on a traveling sound of an engine driven vehicle, which has been generally used in the past, there is a problem in that the pedestrian or the bicycle rider may be slow in noticing presence of an electric motor driven vehicle (or a hybrid automobile or a fuel-cell-powered automobile) with a high silence characteristic or may make an error in assumption of a distance.\nUsually, a pedestrian walking on a street or a bicycle rider can recognize approach or the like of an automobile (a vehicle) according to an engine sound or the like of the vehicle. However, in the case of the hybrid automobile, since the hybrid automobile is in a mode of traveling by an electric motor rather than traveling by an engine at the time of low-speed traveling, no engine sound is emitted, which causes a near collision or the like. Similarly, in the case of the electric automobile or the fuel-cell-powered automobile, the automobile is always in a mode of traveling by an electric motor. Thus, since a pedestrian cannot recognize approach of the electric automobile or the fuel-cell-powered automobile, which causes a near collision as in the case of the hybrid automobile.\nTherefore, vehicle-presence notifying apparatuses for solving the problems caused by the silence characteristic provided in the hybrid automobiles, the fuel-cell-powered automobiles, the electric automobiles, and the like have been disclosed in patent documents.\nConcerning this type of a conventional technique, for example, Japanese Patent Application Laid-open No. H11-27810 discloses a technique for an electric automobile that gives an alarm by voice using a speaker or the like to a pedestrian and a driver around its vehicle when the vehicle detects a human body (a pedestrian) with a photographing camera or a radar provided in the vehicle. Japanese Patent Application Laid-open No. H7-322403 discloses a technique for notifying a driving state of an electric automobile to a pedestrian or a driver with sound. Japanese Patent Application Laid-open No. H10-83187 discloses a technique for a hybrid automobile that emits different artificial sounds according to a driving mode at the time of traveling. Japanese Patent Application Laid-open No. H10-178704 discloses a technique for a hybrid automobile that reduces an engine sound When a road, on which a vehicle is traveling, is congested.\nHowever, vehicle-presence notifying apparatuses according to the conventional techniques have problems as described below. In the case of the conventional technique disclosed in the Japanese Patent Application Laid-open No. H11-27810, it is necessary to provide a sound source for the speaker or the like anew in the electric automobile, which causes an increase in cost. In the case of the conventional techniques disclosed in Japanese Patent Application Laid-open No. H7-322403 and Japanese Patent Application Laid-open No. H10-83187, an artificial sound for notifying a driving state of the vehicle to a pedestrian or a driver is simply emitted or an artificial sound is emitted according to a driving mode, which only prevents a silence characteristic originally provided in the hybrid automobile and the fuel-cell-powered automobile. In addition, in the case of such conventional techniques, it is possible to notify approach of a vehicle to a pedestrian with an artificial sound or the like emitted from the speaker. However, in the notifying by such an artificial sound, since a pedestrian not in a traveling direction of the vehicle can also hear an alarm sound, the technologies are contrary to the silence characteristic originally provided in the hybrid automobile and the fuel-cell-powered automobile. In residential areas in the suburbs rather than a city where there is heavy traffic of pedestrians, noise caused by automobiles is a problem. Thus, notifying by an alarm sound using the speaker or the like is not so preferable in terms of the environment."}
{"text": "1. Field of the Invention\nThe invention relates to a method of manufacturing a metal foil for an electrolytic capacitor, which is used for an anode foil or a cathode foil, and an electrolytic capacitor using the same.\n2. Description of the Related Art\nIn the modern society, automobiles and electric appliances such as televisions and personal computers have been so widely spread as to be indispensable for daily lives. High performance of electronic parts such as electrolytic capacitors greatly contributes to high functionalization of these products.\nIn a general electrolytic capacitor, a metal foil made of a valve metal such as aluminum, tantalum, titanium and the like is used for an electrode foil. For example, an electrolytic capacitor, produced by forming a capacitor element in which a rectangular anode foil and a rectangular cathode foil are wound, impregnating a driving electrolytic solution between the rectangular anode foil and the rectangular cathode foil and housing the capacitor element in a casing, has been commercialized. Also, a solid electrolytic capacitor containing a solid electrolyte excellent in electric conductivity such as polythiophene and polypyrrole, instead of the driving electrolytic solution, has been commercialized.\nA metal foil for use in such electrolytic capacitors, for example, is manufactured by following method. First, a wide metal strip is subjected to etching treatment to widen the surface area. Next, a wide metal strip for an anode foil having an anodic oxide film which is formed by anodizing the surface of the etched wide metal strip and a wide metal strip for a cathode foil which is only performed the etching treatment are prepared. Thereafter, in a slitting step, these respective wide metal strips are slit in predetermined width with a cutting blade. Further, in a cutting step, the slit metal sheets are cut in pieces by a cutting apparatus to produce rectangular metal foils. (See, for example, Japanese Unexamined Patent Publication No. 2007-152436)\nRecently, a variety of electronic parts tend to be installed in automobiles and the like for high functionalization. Therefore, electrolytic capacitors with high withstanding voltage having a rated voltage as high as 25 V or 35 V have been required.\nHowever, such electrolytic capacitors with high withstanding voltage have a problem that a short circuit is likely to be caused at the time of manufacture. According to investigations of the present inventor, it is supposed that burrs and cracks produced at a slit edge surface of an anode foil or a cathode foil to be used for manufacturing an electrolytic capacitor with high withstanding voltage are causes of the problem. Specifically, in the case where the anode foil and the cathode foil of an electrolytic capacitor are produced, as described above, by carrying out the slitting step of slitting the wide metal strip subjected to etching treatment in predetermined width, it is required to use metal foils with large etching pits as the anode foil and the cathode foil to produce the electrolytic capacitor with high withstanding voltage. Therefore, when the wide metal strip is slit, the unevenness of the slit edge surface becomes large and accordingly, burrs and cracks tend to be formed at the slit edge surface. Particularly, because the hard anodic oxide film is formed on the surface of the wide metal strip for the anode foil, the cutting blade becomes dull and whereby, burrs and cracks are easy to be formed at the slit edge surface.\nIn the case where the metal foil having burrs and cracks at the slit edge surface is used to manufacture an electrolytic capacitor, both electrode foils are electrically conducted with the burrs and whereby, a leak current tends to be increased or a short circuit tends to be occur. Further, due to progression of cracks, the leak current tends to be further increased and the capacity tends to be lowered. Therefore, in the industrial manufacture, there is a problem that defective products causing the short circuit and the like tend to be increased and the production yield tends to be lowered."}
{"text": "The present invention relates generally to satellite maneuvering systems, and more particularly, to a magnetic dipole tractor beam control system.\nSpacecraft are directed through various maneuvers to guide them along a variety of celestial paths (e.g., interplanetary voyages and orbits around astronomical bodies such as the Earth) to perform a variety of functions (e.g., weather or planet surface monitoring, commercial communications and scientific observations). These spacecraft maneuvers include ones which initially place a spacecraft in a predetermined celestial path, ones which maintain a desired spacecraft station in the celestial path and ones which maintain a desired spacecraft attitude in the celestial path. In an exemplary case in which the celestial path is a geosynchronous Earth orbit, these spacecraft maneuvers are typically referred to as orbital transferring, stationkeeping and attitude controlling.\nSpacecraft maneuvers are accomplished by the application of forces and torques to the spacecraft. Typically, orbital transferring and stationkeeeping are achieved by directing force vectors through the center of mass of spacecraft so as to obtain changes in spacecraft position without disturbing spacecraft attitude. In contrast, attitude controlling is typically achieved by spacing force vectors from the center of mass to thereby generate torque vectors which realize attitude changes.\nConventional spacecraft structures for application of spacecraft forces and torques include thrusters, momentum and reaction wheels, surfaces which receive solar pressure (e.g., solar cell arrays), extended masses which interact with ambient gravity gradients, mechanical inter-spacecraft control structures (e.g., mechanical arms) and magnetic torquing coils which interact with ambient magnetic fields (e.g., the Earth\"\"s magnetic field).\nAlthough all of these control structures have been used to effect spacecraft maneuvers, each has characteristics that limit their usefulness. Thruster systems are typically bulky, heavy and expel propellant products which can coat and degrade sensitive spacecraft structures (e.g., electro-optical instruments and solar cell arrays). In addition, fuel is an expendable substance of limited supply and, hence, its lack routinely produces the effective end of useful spacecraft life. Momentum and reaction wheels are restricted to the application of torques and their momentum must be periodically xe2x80x9cdumpedxe2x80x9d with other control structures (e.g., limited-fuel thrusters) when it approaches the design limit of the wheels.\nGeneration of spacecraft forces and spacecraft torques by use of solar pressure and ambient gravity gradients typically requires the arrangement or deployment of mechanical structures (e.g., selective rotation of solar cell arrays or extension of gravitational masses from the spacecraft on booms or tethers). Alternatively, generation of spacecraft forces and spacecraft torques can be effected with the limited quantity of thruster fuel.\nInter-spacecraft mechanical control structures are typically bulky and heavy which means they use a significant amount of spacecraft volume and weight (quantities which are always in short supply). Such structures can only be applied when the spacecraft spacing is less than the maximum reach of the control structures. In addition, direct-contact mechanical control may initiate a damaging electrostatic discharge because of spacecraft potential differences, may cause mechanical damage and may respond to electrical or mechanical failure by failing to uncouple the spacecrafts.\nThe application of magnetic fields to spacecraft maneuvers has typically been directed to the use of ambient magnetic fields (e.g., the Earth\"\"s magnetic field) or to the theoretical use of magnetic structures that have been previously distributed in orbit about a celestial body (e.g., see Lebon, Benoit A., xe2x80x9cMagnetic Propulsion along an Orbiting Grain Streamxe2x80x9d, Journal of Spacecraft and Rockets, Vol. 23, March-April, 1986, pp. 141-143).\nThe forces and torques generated in magnetic structures by an ambient magnetic field have been generally described by many investigators (e.g., see Boyer, Timothy H., xe2x80x9cThe Force on a Magnetic Dipolexe2x80x9d, American Journal of Physics, August 1988, Vol. 56, No. 8, pp. 688-692; Brownstein, K. R., xe2x80x9cForce Exerted on a Magnetic Dipolexe2x80x9d, American Journal of Physics, October 1993, Vol. 61, No. 10, pp. 940-941); Greene, Jack B., et al., xe2x80x9cForce on a Magnetic Dipolexe2x80x9d, American Journal of Physics, February 1971, Vol. 39, pp. 172-175; Hnizdo, V., xe2x80x9cHidden Momentum and the Force on a Magnetic Dipolexe2x80x9d, Magnetic and Electrical Separation, 1992, Vol. 3, pp. 259-265; and Vaidman, Lev, xe2x80x9cTorque and Force on a Magnetic Dipolexe2x80x9d, American Journal of Physics, October 1990, Vol. 58, No. 10, pp. 978-983).\nMagnetic forces on neutrally charged objects are not induced by uniform magnetic fields. Accordingly, ambient magnetic fields cannot be used to generate forces on spacecraft because they are essentially uniform at the spatial scale of spacecraft. In addition, the generation of torques with an ambient magnetic field is limited in application because the direction of the ambient magnetic field gradient cannot be selected.\nU.S. Pat. No. 6,089,510 teaches using a magnetic dipole for spacecraft maneuvers. In the \"\"510 patent, spacecraft maneuvers of a first spacecraft are realized with conventional force and torque generators. Spacecraft maneuvers of a second spacecraft are realized through magnetic interaction between the first and second spacecraft. In particular, a magnetic moment vector m1 of a magnetic system of the first spacecraft and a magnetic moment vector m2 of a magnetic system of the second spacecraft are adjusted to apply selected force vectors and torque vectors to the first and second spacecraft. Unfortunately, the system described in the \"\"510 patent uses open loop control.\nThe disadvantages associated with these conventional magnetic maneuvering techniques have made it apparent that a new technique for control of a magnetic dipole tractor beam is needed. Preferably, the new technique would provide closed loop control of a magnetic dipole tractor beam. The present invention is directed to these ends.\nIt is, therefore, an object of this invention to provide an improved and reliable magnetic dipole tractor beam control system. Another object of the invention is to provide closed loop control of a magnetic dipole tractor beam.\nIn accordance with the objects of this invention, a magnetic dipole tractor beam control system is provided. In one embodiment of the invention, spacecraft maneuvers (e.g., orbit transferring, stationkeeping and attitude controlling) of a first spacecraft are realized with conventional force and torque generators (e.g., thrusters and momentum wheels). Spacecraft maneuvers of a second spacecraft are realized through magnetic interaction between the first and second spacecraft using a closed loop control system. In particular, a magnetic moment vector m1 of a magnetic system of the first spacecraft and a magnetic moment vector m2 of a magnetic system of the second spacecraft are adjusted to apply selected force vectors and torque vectors to the first and second spacecraft using a closed loop control system.\nThe present invention thus achieves an improved magnetic dipole tractor beam control system. The present invention is advantageous in that it is capable of providing closed loop control of a magnetic dipole tractor beam.\nAdditional advantages and features of the present invention will become apparent from the description that follows and may be realized by means of the instrumentalities and combinations particularly pointed out in the appended claims taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates in general to semiconductor memories and in particular to a non-volatile memory array architecture and a method of operation of the same.\nThe most common variety of non-volatile memories, such as EPROM, flash memory, and some EEPROMs, today employs channel hot electron (CHE) for programming and negative gated Fowler-Nordheim (FN) tunneling for erase. FIG. 1 shows a conventional n-channel stack gate flash memory cell 100. Memory cell 100 includes N+ source 102 and drain 103 regions spaced apart in a P-type silicon substrate 101 to form a channel region 104 therebetween. A floating gate 105 of polycrystalline silicon material is laid on top of a tunneling dielectric 106, which extends over the channel region 104 and overlaps the source 102 and drain 103 regions. Stacked on top of, but insulated from, floating gate 105 is a gate 107 of polycrystalline material. Junction 102 is made deeper than normal in order to minimize the adverse reliability effects of tunnel oxide hot hole trapping during erase operation.\nCell 100 is programmed, i.e., its threshold voltage is raised higher, by applying 10V to gate 107, 5V to drain 103, and grounding source 102. The memory cell is thus strongly turned on, and the cell\"\"s threshold voltage is raised due to injection of hot electrons from the channel region near the drain 103 to floating gate 106, as indicated by the arrow labeled as xe2x80x9cPxe2x80x9d. Cell 100 is erased, i.e., its threshold voltage is lowered, by applying xe2x88x9210V to gate 107, 5V to source 102, and floating drain 103. The cell\"\"s threshold voltage is thus lowered due to tunneling of electrons from the floating gate 105 to source 102, as indicated by the arrow labeled as xe2x80x9cExe2x80x9d.\nConventional memory arrays include a matrix of memory cells arranged along rows and columns. The gates of the cells along each row are connected together forming a wordline. In one array architecture, the cells along each column are grouped in a number of segments, and the drains of the cells in each segment are coupled to a corresponding segment line. The segment lines along each column are coupled to a corresponding data line through one or more segment select transistors. The segmentation of the cells in each column helps reduce the bitline capacitance to that of the metal bitline plus the small capacitance of a selected segment line. Performance of the memory is thus improved.\nDuring programming or read operations, one or more bitlines are selected through a column select circuit for transferring data to or from the selected memory cells. The column select circuit typically has a multiplexer configuration in that a group of serially-connected NMOS pass transistors controlled by column decoding signals selectively couple one or more bitlines to either sense amplifiers (read operation) or data-in buffers (programming operation). Depending on the total number of bitlines in the array and the number of bitlines to be selected, two or more levels of column selection need to be implemented in the column select circuit. The number of levels of column selection correspond to the number of serially-connected pass transistors that couple the selected bitlines to the sense amplifier or data-in buffer. For example, if two levels of column selection are implemented, a selected bitline will be coupled to the sense amplifier or data-in buffer through two serially-connected column select transistors.\nThe sizes of the column select transistors and the segment select transistors need to be made large enough so that the required cell programming voltage and current can be provided to the selected cell. Because of the programming biasing conditions, the serially connected segment select transistor and column select transistors result in a rather resistive path, which can be compensated for by increasing the transistor sizes. This can be more clearly understood with the help of FIG. 2.\nFIG. 2 shows a portion 201 of an array along with a portion 202 of a column select circuit. The array portion 201 includes a memory cell 100 with its gate coupled to wordline WL and its drain coupled to a segment line S0. The source of cell 100 is shown as being connected to ground for simplicity, although, the source is typically connected to a source line which may be decoded to provide ground only to selected memory cells. A segment select transistor MS is coupled between segment line S0 and bitline BL, with its gate coupled to segment select signal SS. Bit line BL is coupled to the data-in block 204 through two serially connected column select transistors MYa and MYb. Column select transistors MYa and MYb are controlled by column decode signals Ya and Yb, respectively. As indicated in FIG. 2, the deeper source junction 102 of the FIG. 1 cell is connected to ground, while the shallower drain junction 103 is connected to segment line S0.\nAs can be seen, cell 100, and transistors MS, MYa, and MYb are serially-connected to data-in block 204. To program cell 100, 10V is supplied to wordline WL, while 5V needs to be supplied to its drain, i.e., to segment line S0. To supply 5V to segment line S0, data-in block 204 outputs 5V on line 206, and column select signals Ya and Yb as well as segment select signal SS are raised to 10V. Thus, the 5V on line 206 is transferred through transistors MYa, MYb, and MS to segment line S0. The drive capability of each transistor MYa, MYb, MS is, to a first order approximation, equal to its Vgsxe2x88x92Vt, wherein Vgs represents the transistor gate to source voltage, and Vt represents the transistor threshold voltage. For each of transistors MYa, MYb, MS, Vgs=Vgxe2x88x92Vs=10Vxe2x88x925V=5V, and the Vt is approximately 2V because of the back bias effect. Thus, for each transistor MYa, MYb, MS, Vgsxe2x88x92Vt=5Vxe2x88x922V=3V. Because of the small Vgsxe2x88x92Vt of 3V, the sizes of these transistors need to be made large so that sufficient current can be supplied to the cell during programming.\nIn higher density memories, where the number of levels in the column select circuit increases, the sizes of the column select transistors increase proportionally. This increases the die size. More importantly, as higher performance is required of memory devices, the need for further segmentation of the bitlines increases, resulting in a larger number of segment select transistors in the array. The adverse impact of a larger size segment select transistor and a larger number of segment select transistors on the overall die size can be rather substantial.\nFIG. 3 illustrates another draw back of conventional memory arrays, namely, the non-uniform programming characteristic of memory cells in the array due to the source resistance. A portion 300 of a memory array is shown as including 16 memory cells 100-0 to 100-15 along a row. The drain of each cell is coupled to a corresponding segment line S0 to S15, and the gates of the cells are connected to a wordline WL. The sources of the cells are connected together and to metal source lines SLn and SLn+1 through a diffusion strip 310. Resistors R0 to R16 depict the resistance associated with the diffusion strip 310. The cell configuration of FIG. 3 is repeated as many times as required to form the entire array.\nFor the above-indicated cell biasing during programming, the cell programming performance is dependent primarily upon the gate to source voltage Vgs of the cell. For example, with the wordline WL at 10V, and the source fully grounded, the cell Vgs equals a full 10V. However, because of the presence of the resistive diffusion strip 310, depending on the location of the cells along the diffusion strip 310, the effective Vgs of the cells vary. For example, of the 16 cells, the cells closest to the center of the diffusion strip will have the maximum source resistance, and thus poorer programming characteristics, while the cells closest to the ends of the diffusion strip 310 have minimum source resistance, and thus the best programming characteristics. This leads to the undesirable non-uniform programming characteristics of the cells across the array.\nThus, an array architecture and method of operation are needed whereby the adverse effect of column select and segment select transistor sizes on the die size can be minimized, while a more uniform programming characteristic across the array cells can be obtained.\nIn accordance with an embodiment of the present invention, an array of non-volatile memory cells are arranged along rows and columns. Each memory cell has a drain region spaced apart from a source region to form a channel region therebetween. The drain region has a greater depth than the source region. Each memory cell further has a stack of floating gate and select gate extending over the channel region, the select gate of the cells along each row being connected together to form a wordline. Each of a plurality of data lines is coupled to the drain region of at least a portion of a column of cells. Each of a plurality of source lines is coupled to a source region of a plurality of cells along at least a portion of a row of cells, wherein injection of hot electrons from a portion of the channel region near the source region to the floating gate is induced in a selected memory cell in the array by applying a first voltage to a selected data line to which the drain of the selected memory cell is coupled, a second positive voltage to a word line to which the selected gate of the selected memory cell is coupled, and a third positive voltage to a source line to which the source of the selected memory cell is coupled, wherein the injection of hot electrons increases a threshold voltage of the selected cell.\nIn one embodiment, the first voltage is in the range of xe2x88x920.5 to +0.5V, the second positive voltage is in the range of 8V to 12V, and the third positive voltage is in the range of 4.0V to 6.0V.\nIn accordance with another embodiment of the present invention, a method of operating a memory array is provided. The memory array has a plurality of memory cells arranged along rows and columns. Each memory cell has a source region spaced apart from a drain region to form a channel region therebetween, and a floating gate and select gate stack extending over the channel region. The select gates of the memory cells along each row are coupled together to form a plurality of wordlines. Each of a plurality of data lines is coupled to a drain region of each of a plurality of non-volatile memory cells along a column. Each of a plurality of source lines is coupled to a source region of each of a plurality of memory cells along a row. The method of operation of the memory array is as follows. A first voltage is applied to a selected data line coupled to the selected cell. A second positive voltage is applied to a selected word line coupled to the selected cell. A third positive voltage is applied to a source line coupled to the selected cell, wherein upon applying the first voltage and the second and third positive voltages an injection of hot electrons from a portion of the selected cell\"\"s channel region substantially near the source region to the selected cell\"\"s floating gate is induced whereby a threshold voltage of the selected memory cell is increased.\nIn one embodiment, the first voltage is in the range of xe2x88x920.5 to +0.5V, the second positive voltage is in the range of 8V to 12V, and the third positive voltage is in the range of 4.0V to 6.0V.\nOther advantages and features of the present invention will be apparent from the accompanying drawings and the description thereof."}
{"text": "The advent of hearing aids which are inserted directly into the user's ear canal has highlighted the problem of ear wax (cerumen) build up. Cerumen produced in the ear canal can damage hearing aids in several ways. Firstly, the wax can build up in the sound spout of the hearing aid and reduce the sound output by a significant amount. When cerumen totally blocks the sound spout the sound output is nil and the hearing aid appears dead.\nSecondly, cerumen may migrate through the sound spout and into the hearing aid receiver thereby causing irreparable damage to the transducer and necessitating extensive repair work.\nIn addition, wax accumulation presents a problem even if the sound spout is only partly blocked because removing the wax from the hearing aid during cleaning may result in wax inadvertently falling down the sound spout into the receiver port.\nAs a result of the wax build-up and migration problem, various attempts have been made to delay the progress of the wax into the sound spout and receiver port of the hearing aid. Such attempts are shown, for example, in U.S. Pat. No. 4,553,627 filed Oct. 22, 1984, European Patent Application No. 0287315 filed Dec. 4, 1988, and U.K. Patent Application GB2155276A filed Jun. 18, 1984. The above patents show various approaches which have been taken to provide a wax guard for the sound spout of a hearing aid. However these proposals have had only limited success in delaying the progress of wax into the sound spout and receiver port or the hearing aid. The ease of removal and cleaning of the guard is very important given the advanced age of many hearing aid users."}
{"text": "Industrial computer systems commonly used for process control purposes are typically made up by combining a plurality of system components made by different manufacturers to meet the performance requirements of different system applications. While it is frequently possible to adapt each individual component for use in a harsh factory environment by providing sealing means in each compartment for excluding dust, moisture, chemical fumes and the like, the provision of such sealing means on an individual component basis is expensive and typically results in components which are characterized by the difficulties encountered in servicing the components. More important, the need for providing system components which are individually sealed seriously limits the variety of systems which can be set up using commercially available components. It has been found that the provision of sealed compartments for such systems has resulted in similar difficulties with respect to permitting free use of various components in the harsh factory environments while also permitting easy servicing of system components in a system console which facilitates operation of a system by a system operator.\nIn an effort to resolve these problems, factory floor operator interface consoles have been proposed, such as that disclosed in U.S. Pat. No. 4,827,439. While this console provides a sealed compartment chamber for a plurality of process control or computer system components including a central processing unit of a computer, a video display unit having a screen at one end for displaying data provided by the computer, and typically a recorder unit or the like, it does not provide protection for the keyboard when not in use, or a sealed enclosure for a printer, or means to access printer output from a position external to the sealed enclosure.\nU.S. Pat. No. 4,861,121 discloses a space-efficient cabinet for housing a computer workstation. This patent discloses a space-efficient cabinet for housing a computer workstation which can include a video display monitor, a central processing unit with disk drive and a keypad. This cabinet has two generally opposed sidewalls and generally opposed front and back walls, which extend between the sidewalls. At least one horizontal surface is provided extending between the front, back and opposed sidewalls to carry the video display monitor. An inner angled shelf is also present, and is affixed between the sidewalls, to carry the central processing unit and disk drive. The shelf is angled downwardly from front to back relative to the horizontal allowing access to the central processing unit and disk drive whereby disks can be readily inserted into or removed from the central processing unit and disk drive. The angle of the shelf carrying the central processing unit and disk drive ensures a minimal footprint of the cabinet on the floor thus saving floor space. A movable partition is also provided to support the keypad. However, this device does not provide for any sealed environment to protect the computer system from a harsh industrial environment, and does not provide a printer enclosure or means to have access to the printer output from externally of the cabinet.\nU.S. Pat. No. 4,718,740 discloses a housing and storage mechanism for a terminal keyboard and display panel. This patent discloses a rectangular housing mounted vertically on a wall having a front face hinged along the lower edge for opening to a horizontal position. A keyboard is mounted on the upper portion of the housing face interiorly of the housing. A data display panel is pivotally suspended from the lower edge of the keyboard. A mechanism activated by the front face engages the lower edge of the display panel and elevates the panel to an angle inclined for convenient viewing when the front face is lowered to a horizontal position. This device is similar to U.S. Pat. No. 4,704,604 which provides for the secure positioning of the keyboard adjacent to the video display unit in either a lowered, use position or a raised, stored position. In the use position, the keyboard is maintained in a stable horizontal orientation immediately below and forward of the face plates of the video display unit. In the stored position, the keyboard is securely maintained in a generally vertical orientation immediately forward of the video display unit's face plate so as to provide protection therefore and enhance the portability of the computer terminal. While these devices provide for the protection of the keyboard and the video display terminal, there is no provision for enclosing and protecting the central processing unit and a printer, or for providing access to printer output external to the enclosure.\nU.S. Pat. No. 4,930,917 discloses construction for a sound reducing printer box. However, this soundproof box construction is not designed to protect the printer from damage in an industrial environment, and has no provision for enclosing the other components of the computer system, or for providing access to printer output external to the enclosure."}
{"text": "The search for novel anti-cancer drugs is a never ending story, since cancer is becoming a more and more important cause of death amongst humans. More than 80% of all anti-cancer drugs are directed towards the apoptosis pathway of tumour cells and are cytotoxic upon activating said pathway.\nA large number of cancer cells such as glioblastomas (brain cancers), brain metastases, melanomas, pancreatic cancers, lung cancers of the NSCLC-type, refractory prostate cancers (HRPC), breast cancers such as triple negatives and other types are naturally resistant to apoptosis and cannot be treated by the many known drugs and chemotherapeutics. Development of new compounds having cytotoxic and/or cytostatic effects in tumour or cancer cells, and particularly in tumour or cancer cells which are resistant to apoptosis, is of considerable importance.\nA naturally-occurring beta carboline, harmine, has been previously isolated from plants such as the Middle Eastern plant harmal or Syrian rue (Peganum harmala) and the South American vine (Banisteriopsis caapi) and shown to exhibit anti-cancer properties. However, previous reports demonstrated that harmine and its derivatives caused remarkable acute neurotoxicity characterized by tremble, twitch, and jumping in experimental mice model. Results of investigation on the in vitro anti-tumour activity of harmine and certain of its derivatives showed that these compounds had significant inhibition effect on several cultured tumour cell lines, such as HeLa cells (cervical carcinoma), S-180 cells (sarcoma), BEL-7402 cells (hepatoma), MGC-803 cells (gastric carcinoma), CNE2 cells (nasopharyngeal carcinoma), MA782/5S cells (breast cancer) and K562 cells (leucocythemia) (Wu, US 20090227619). EP 0357122 showed cytostatic properties of certain beta-carbolines.\nMoreover, harmine and certain of its derivatives exhibited DNA intercalation capacity and topoisomerase I inhibition, which was related to their anti-tumour activity (Cao et al., 2005, Biochem. Biophys. Res. Com. 338: 1557-63).\nAdditionally, harmine was repeatedly shown to inhibit the protein kinase DYRK1A (Göckler, 2009, FEBS J. 276(21): 6324-37; Seifert, 2008, FEBS J. 275(24): 6268-80; Bain et al., 2007, Biochem J 408: 297-315). Over-expression of the DYRK1A kinase (dual specificity tyrosine-phosphorylated and regulated kinase 1a) has been implicated in multiple diseases or disorders or syndromes prominently including cancer, tumorigenesis and uncontrolled cell proliferation (Laguna et al., 2008, Dev Cell. 15(6): 841-53), but also Alzheimer disease, Pick disease (Ferrer et al., 2005, Neurobiol Dis. 20(2): 392-400) and Down Syndrome (Lepagnol-Bestel et al., 2009, Hum Mol Genet. 18(8): 1405-14).\nMoreover, caspase-9 is a substrate of DYRK1A: inhibition of DYRK1A, by, e.g., harmine, induces the activation of caspase-9 and leads to massive apoptosis in different human cell types (Seifert et al., supra). It was also shown that melanomas are intrinsically resistant to pro-apoptosis stimuli and over-express DYRK1A (de Wit et al., 2002, Melanoma Res 12: 57-69). Recently, it was shown that harmine could reverse resistance to anticancer drugs by inhibiting BCRP (Ma et al., 2010, Phytother. Res. 24(1): 146-9).\nThere is therefore a continuous need in the art for identifying novel anti-proliferative drugs and for improving the efficacy of existing anti-proliferative treatments by combining the latter with novel anti-proliferative drugs."}
{"text": "More recently, drilling tools as well as milling tools have been developed, which, contrary to integral solid tools, are composed of two parts, viz. a basic body and a head detachably connected with the same and thereby being replaceable, in which head the requisite cutting edges are included. In such a way, the major part of the tool can be manufactured from a comparatively inexpensive material having a moderate modulus of elasticity, such as steel, while a smaller part, viz. the head, can be manufactured from a harder and more expensive material, such as cemented carbide, cermet, ceramics and the like, which gives the cutting edges a good chip-removing capacity, good machining precision and long service life. In other words, the head forms a wear part, which can be discarded after wear-out, while the basic body can be re-used several times. A now recognized denomination of such, cutting edge-carrying heads is “loose tops”, which henceforth will be used herein together with the concept “loose top tools”.\nOn rotatable tools of the loose top type, several requirements are made, one of which is that the loose top should be held centered in an exact and reliable way in relation to the basic body. Accordingly, each unintentional eccentricity between the center axis of the loose top and the center axis of the basic body should not be more than 0.01 mm. Another requirement or desire from the users' side is that the loose top should be mountable and dismountable in a rapid and convenient way without risk of incorrect mounting. Most preferably, mounting and dismounting should be possible to be carried out without the basic body necessarily having to be removed from the driving machine.\nDrilling tools as well as milling tools (shank-end mills) of the loose top type are widely described in the patent literature and may be divided into a number of different categories depending on the ideas on which the designs are based. Thus, certain tools use loose tops having rear pins that entirely or partly (together with other coupling details) fulfill the task of centering the loose top in relation to the basic body. To this category belongs among others the tool disclosed in U.S. Pat. No. 6,012,881, which discloses a loose top drill in which a rear coupling part of the loose top is axially inserted in a jaw between two non-compliant drivers, the insides of which include axially extending, torque-transferring ridges that engage the corresponding chutes in the coupling part included in the loose top, besides which a centric pin protruding rearward from the coupling part is inserted in a center hole mouthing in the bottom of the jaw. With the pin, a screw mounted in a radial, threaded hole in the basic body co-operates, which has the purpose of locking the loose top in relation to the basic body. In that connection, the center pin is cylindrical and insertable at a close (form fitting) fit in a likewise cylindrical center hole, in order to, together with concave and convex contact surfaces of the drivers and the coupling part, respectively, fulfill the purpose of centering the loose top. However, the concurrent requirements of fit not only between the center pin and the hole, but also between the drivers and the coupling part, impose extreme, not to say unattainable, requirements of manufacturing precision. In addition, a conflicting problem will arise if good precision peradventure would be achieved, viz. that the mounting and dismounting of the loose top become difficult to carry out as a consequence of the centering pin requiring great force to be pushed into and pulled out of, respectively, the hole.\nWith the purpose of obviating the above-mentioned disadvantages of the tool according to U.S. Pat. No. 6,012,881, it has by EP 2266736 and EP 2266734 been proposed to form the centering pin of the loose top with, on one hand, an essentially semi-cylindrical contact surface concentric with the center axis of the loose top, and on the other hand a diametrically opposite clearance surface, the first-mentioned one of which can—by means of the screw of the tool—be pressed against the inside of a cylindrical centering hole in the basic body. In such a way, the centering pin can be given a cross-sectional area that is smaller than the cross-sectional area of the hole, the mounting and dismounting of the loose top being facilitated without the exact centering of the loose top being jeopardized. However, in this case, a 50% risk arises of incorrect mounting of the loose top in connection with the centering pin being inserted into the hole, in that the contact surface of the pin can be turned toward any half of the internal, cylindrical hole wall.\nThe present invention aims at obviating not only the serious disadvantages inherent in the loose top tool according to U.S. Pat. No. 6,012,881, but also the practical mounting inconveniences that have turned out to be associated with the tools and loose tops that are objects of EP 2266736 and EP 2266734.\nAn object of the invention to provide a loose top tool the loose top of which can, on one hand, be centered in an accurate way in relation to the basic body of the tool, and on the other hand be mounted in a simple way without risk of incorrect mounting. In other words, the operator should without mental effort be able to, in a foolproof way, mount the loose top in only one predetermined position.\nAnother object of the invention is to provide a loose top tool, the two main components of which, i.e., the basic body and the loose top, respectively, should be manufacturable in an effective and economical way while achieving the desired simplicity of the mounting.\nYet another object of the invention is to provide a loose top the centering pin of which is strong and robust."}
{"text": "This invention relates to an improved oxidant and more in particular to an acidic, aqueous, oxidizing agent containing bromate and iodate ions.\nDying of various fabrics to impart a color to the fiber has been practiced for many centuries. The color must generally be permanently and uniformly distributed throughout the fiber and not merely superficially applied to the fiber as in painting. Many different types of natural and regenerated cellulosic fibers have been dyed to impart a color. For example, natural fibers, such as the vegetable fibers cotton, linen, jute, and flax have been dyed. Regenerated cellulosic fibers, such as viscose rayon and cellulose acetate, are those produced from natural materials which were altered by man to produce a desired textile material.\nIt has become accepted, and common, practice to color these materials with well-known sulfur and vat dyes. These dyes are water insoluble substances which are readily converted to a water soluble or leuco form by reducing the sulfur or vat dye in, for example, a solution containing an alkali and sodium sulfide or hydrosulfite.\nThe leuco forms of sulfur and vat dyes are water soluble and well known to be substantive to cellulosic fibers. After application to the fiber, the leuco dye must be oxidized to permanently color the fabric. The process of U.S. Pat. No. 3,775,047 oxidized the dye with an aqueous oxidizing solution including acetic acid and sodium or potassium iodate. U.S. Pat. No. 4,042,319 disclosed similar oxidation with an aqueous oxidant containing acetic or formic acid, an alkali bromate and an alkali iodate. Such oxidizing solutions are operable; however, it is desired to provide an improved material suitable to oxidize leuco forms of sulfur and vat dyes."}
{"text": "Subterranean coal seams contain substantial quantities of natural gas, primarily in the form of methane. Such reservoirs are commonly referred to as coal-bed methane (CBM) reservoirs. To effectively produce gas from CBM reservoirs, one or more stimulation methods are used to increase the output.\nThe most common CBM stimulation methods include directionally drilling parallel to the bedding of the seam, the creation of a cavity in a coal seam, or hydraulic fracturing of the coal.\nDirectional drilling involves angling the drill stem so that drilling is not vertical but will parallel the coal seam. Because the bore hole trends along the formation, a greater area of the wellbore is in contact with the coal seam and thus higher gas extraction is possible.\nThe cavitation completion method creates on or more cavities in the coal. The purpose is to reduce the amount of damage to the surrounding structure that may have resulted during drilling, to create an enhanced permeability zone outside of the cavitated area, and to reduce the near wellbore flow resistance arising from the convergence of flow in a radially inward flow field. Typically cavity completions are performed on open-hole completed wells (no casing across the production interval) in a cyclic manner thus the treatment is usually referred to as cavitation cycling as described for example in: Palmer, I. D., Mavor, M. J., Spitler, J. L., Seidle, J. P., and Volz, R. F. 1993b. Openhole cavity completions in coalbed methane wells in the San Juan Basin. Journal of Petroleum Technology, 45(11):1072-1080 (November).\nCompressed air is most frequently used to pressurize the near wellbore zone. These gases can be foamed to reduce their apparent mobility in the cleat network and contain them in the near wellbore region where they will be most effective. The pressure is bled off as fast as possible, rubblizing the near wellbore zone and producing it out of the wellbore, creating a rather large cavity surrounding the wellbore. The procedure reduces the wellbore skin effect.\nCavitation cycling uses several mechanisms to link the wellbore to the coal fracture system. These mechanisms include creating a physical cavity in the coals of the open-hole section (up to 10 feet in diameter); propagating a self-propping, vertical, tensile fracture that extends up to 200 feet away from the wellbore (parallel to the direction of maximum stress and perpendicular to the minimum principal stress); and creating a zone of shear stress-failure that enhances permeability in a direction perpendicular to the direction of least stress as described for example in: Khodaverian, M. and McLennan. 1993. Cavity completions: a study of mechanisms and applicability. Proceedings of the 1993 International Coalbed Methane Symposium (Univ. of Alabama/Tuscaloosa), pp. 89-97.\nCavitation is accomplished by applying pressure to the well using compressed air or foam, and then abruptly releasing the pressure. The over-pressured coal zones provide a pressure surge into the wellbore (a “controlled blowout”), and the resulting stress causes dislodgement of coal chips and carries the chips up the well. These cycles of pressure and blowdown are repeated many times over a period of hours or days, and the repeated alternating stress-shear failure in the coal formation creates effects that extend laterally from the wellbore as described for example in: Kahil, A. and Masszi, D. 1984. Cavity stress-relief method to stimulate demethanation boreholes. SPE Paper No. 12843, Proceedings 1984 SPE Unconventional Gas Recovery Symposium (Pittsburg).\nFracturing is another important method to enhance CBM production. As described for example in: Holditch, S. A. 1990. Completion methods in coal seam reservoirs. SPE 20670, Proceedings 65th SPE Annual Technical Conference (New Orleans), p. 533., typically hydraulic fracturing is performed on cased-hole perforated completion wells typically when the coal permeability is less than 20 mD. There are a great many variations of hydraulic fracturing coal formations, but in a rough approximation, the techniques involve injecting a fluid into the formation at sufficient pressure to initiate and propagate a hydraulic fracture, filling the fracture with proppant by continuing injection of a proppant laden fluid, and then flushing the treatment so that the proppant fills the fracture but not the wellbore. Olsen, et al., describe some additional considerations for fracturing coal-bed methane reservoirs in: Olsen, T. N., Brenize, G., and Frenzel, T.: “Improvement Processes for Coalbed Natural Gas Completion and Stimulation,” SPE 84122, presented at the SPE Annual Technical Conference and Exhibition, Denver (Oct. 5-8, 2003).\nConventional hydraulic fracturing technique is described in many literature sources, as well as when applied to coal rocks. Directional drilling cannot be considered as a pure stimulation technique. It is worth to note that both fracturing (either conventional fracturing or cavitation) and directionally drilling simply increases the amount of the coal seam which is in direct contact with the well bore, and, no technological method has been found yet to increase the original porosity of the formation.    Palmer, I. D., Mavor, M. J., Spitler, J. L., Seidle, J. P., and Volz, R. F. 1993b. Openhole cavity completions in coalbed methane wells in the San Juan Basin. Journal of Petroleum Technology, 45(11):1072-1080 (November).    Holditch, S. A. 1990. Completion methods in coal seam reservoirs. SPE 20670, Proceedings 65th SPE Annual Technical Conference (New Orleans), p. 533.    Palmer, I. D., Lambert, S. W., and Spitler, J. L. 1993a Coalbed methane well completions and stimulations. Chapter 14 in AAPG Studies in Geology 38, pp. 303-341.    Olsen, T. N., Brenize, G., and Frenzel, T.: “Improvement Processes for Coalbed Natural Gas Completion and Stimulation,” SPE 84122, presented at the SPE Annual Technical Conference and Exhibition, Denver (Oct. 5-8, 2003).    Khodaverian, M. and McLennan. 1993. Cavity completions: a study of mechanisms and applicability. Proceedings of the 1993 International Coalbed Methane Symposium (Univ. of Alabama/Tuscaloosa), pp. 89-97 or the above paper by Kahil, A. and Masszi, D. 1984.\nFracturing fluid leak-off through the cleat network during the fracturing process is a major limitation of modern hydraulic fracturing methods in CBM reservoirs. Fracturing fluid efficiency is the simple ratio of the volume of the created fracture at the end of pumping divided by the total volume of fluid injected to create the fracture. For obvious reasons, low fracturing fluid efficiency is undesirable. Expense and waste aside, the leaked off fluid can substantially reduce the permeability of the cleat network and defeat the benefit of the stimulation treatment. Moreover, due to various operational constraints (such as maximum injection pressure, maximum injection rate, cost, etc.) low fluid efficiency limits the fracture length that may be achieved.\nFracturing fluids injected into CBM reservoirs create a complex fracture comprised of some dominant channel and numerous minor channels. The minor channels may be part of the cleat network. High leak-off of the fracturing fluid through the cleat network reduces both the rate of growth of the main fracture and the maximum fracture length. Injecting the fracturing slurry at higher flow rates can compensate for leak-off losses, but the leaked off fluid may be trapped in the cleats, blocking the pathway through which reservoir fluids can flow from the reservoir to the wellbore. The cleat aperture can be too small for proppant to enter and prop the cleats in normal practice. Recently, fracturing fluids have been developed that are non-damaging to cleats. Whilst leaking off into the cleat network, they are free of macro-molecules that create immobile plugs in the cleats. Thus, because they are free of polymers and insoluble solids, the non-damaging filtrate is easily displaced from the cleats when the well is put on production. Olsen et al above describe such a fracturing fluid.\nIt is also known to use water-based fluids with dewatering aids and proppant. The presence of water in a porous medium reduces the flow capacity of that medium to other immiscible fluids such as oil or gas. In the case of CBM reservoirs, water from the injected hydraulic fracturing fluids can infiltrate the cleat network and negatively impact the surface properties (mainly wettability) of the coal. These changes can result in reduced dewatering and lead to coal fines migration, which can plug the cleats. U.S. Pat. No. 5,229,017, by Nimerick and Hinkel, describes chemicals, such as butoxylated glycols, that adsorb onto coal surfaces, rendering them hydrophobic. The hydrophobic coal surface maintains the original surface properties of the coal and hinders re-wetting and re-absorption of surfactant present in the fracturing fluid. These surfactants are delivered in the fracturing fluid and result in more rapid coal dewatering and fracturing fluid recovery out of the fracture.\nOther known techniques attempt to control the leak-off into cleats. For example, U.S. Pat. No. 5,474,129 describes a process of using injected gases to perform cavity completions. Water with a foaming agent is added to the gas to create a foam downhole. The foam reduces the rate of gas leaking away from the wellbore through the cleat network. Instead, gas is trapped in the near wellbore region where it adsorbs into the near wellbore coal. Upon depressurization, the trapped gas expands and destroys the coal fabric and promotes the cavity completion process. This patent is referenced because it identifies the role of the cleats in conducting fluids away from a source. Foamed fracturing fluids are often used in CBM reservoirs to control leak-off, to reduce the hydrostatic head in the wellbore at the end of the treatment (and improve fracturing fluid recovery), and to reduce the amount of damaging polymers that enter the cleat network.\nMethods of enhancing coal-bed methane production are considered enhanced recovery techniques. Enhanced recovery techniques involve flooding the CBM reservoir with gases that adsorb to coal more strongly than methane, and thus displace methane or lighter hydrocarbones from the micropore structure of the coal. Several papers describe various aspects of this technique including:    Fulton, P. F., Parente, C. A., Rogers, B. A.: “A Laboratory Investigation of Enhanced Recovery of Methane from Coal by Carbon Dioxide Injection,” SPE/DOE 8930, presented at the 1980 SPE/DOE Symposium on Unconventional Gas Reservoirs, Pittsburg (May 18-21, 1980);    Chaback, J. J., Morgan, D., and Yee, D.: “Sorption Irreversibilities and Mixture Compositional Behavior During Enhanced Coal Bed Methane Recovery Processes,” SPE 35622, presented at the Gas Technology Conference, Calgary (Apr. 28-May 1, 1996);    Zhu, J., Jessen, K., Kovscek, A. R., and Orr, F. M.: “Analytical Theory of Coalbed Methane Recovery by Gas Injection,” SPE Journal, pp 371-379 (December 2003); or    Gorucu, F. B., Jikich, S. A., Bromhal, G. S., Sams, W. N., Ertekin, T., and Smith, D. H.: “Matrix Shrinkage and Swelling Effects on Economics of Enhanced Coalbed Methane Production and CO2 Sequestration in Coal,” SPE 97963, presented at the 2005 SPE Eastern Regional Meeting, Morgantown (Sep. 14-16, 2005).\nThe use of liquid gas fracturing with proppant with mixtures of carbon dioxide and nitrogen or only carbon dioxide is published in general and is used for example in Canada and in the Eastern United States for shallow gas well stimulation. Liquid CO2 or foams/emulsions created by mixtures of CO2 and nitrogen are reported to be undamaging because the fluid will vaporize and be produced after the treatment. Descriptions of the known methods can be found for example in:    Lillies, A. T., and King, S. R.: “Sand Fracturing with Liquid Carbon Dioxide,” SPE 11341, presented at the SPE Production Technology Symposium, Hobbs (Nov. 8-9, 1982);    Yost, A. B., Mazza, R. L., and Gehr, J. B.: “CO2/Sand Fracturing in Devonian Shales,” SPE 26925, presented at the SPE Eastern Regional Meeting, Pittsburg (Nov. 2-4, 1993);    Mazza, R. L.: “Liquid-Free CO2/Sand Stimulations: An Overlooked Technology—Production Update,” SPE 72383, presented at the SPE Eastern Regional Meeting, Canton (Oct. 17-19, 2001); and    Cambell, S. M., Fairchild, N. R., and Arnold, D. L.: “Liquid CO2 and Sand Stimulations in the Lewis Shale, San Juan Basin, New Mexico: A Case Study,” SPE 60317, presented at the SPE Rocky Mountain Regional/Low-Permeability Reservoirs Symposium and Exhibition, Denver (Mar. 12-15, 2000).\nOnce the appropriate bore hole(s) is completed by using one of the above methods, dewatering must occur to reduce the pressure in the formation. Pressure drop in turn promotes methane release from within the coal into the cleats. If the cleats contain a high enough permeability, that is, inter-connectivity, then the methane will flow from the coal into the well bore and can be extracted.\nIn U.S. Pat. No. 6,412,559, there is described a process of stimulating and enhancing methane production in coal reservoirs. The process uses a stronger adsorbing gas (SAG, stronger than methane such as carbon dioxide or H2S), which swells the formation during the fracturing process, the key element of their process is repeated SAG injection and shut-in steps after the stimulation treatment. The shut-in steps may be between 1 day and 1 year in length. The principle reason for the use of a SAG is to promote enhanced coal-bed methane recovery as the SAG preferentially adsorbs onto the coal and displaces adsorbed methane. The fracturing is employed at the beginning of the repeated SAG injection and shut-in steps to improve the infectivity of the SAG, i.e., the fracturing step is to reduce wellbore skin and enable improved CO2 injection.\nIn the published US patent application 20050082058 the reaction of a predetermined gas with the coal is used to induce shrinkage within the coal matrix, thus reducing effective stress and enhancing the fracture void volume. This increased fracture void volume increases coal-bed permeability and resultant increase in methane gas flows. It also allows placement of proppant within the coals to maintain open fractures in the regions surrounding the propped fracture, thus allowing the enhanced fracture system to communicate more effectively with both the natural fracture system and the wellbore and aiding methane recovery.\nIn the U.S. Pat. No. 5,014,788 there is described a method to improve production by introduction of a swelling gas into the coal which after release generates uneven stress fractures. The emphasis of the known method is in the “relatively rapid reduction” in the pressure."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a system and method for easy and rapid instruction of second language skills. More particularly, the present invention relates to a system and method for rapid instruction of second language skills that can accelerate profession-specific language learning skills.\n2. Related Art\nWhile it may be that the English language is the primary, or most used, language in the United States, the number of people in the U.S. who do not speak English is increasing daily. For instance, it has been estimated that the number of U.S. residents who speak Spanish may shortly outnumber the residents who speak English. This increase of non-English speaking residents posses difficulties for many businesses and organizations. For example, many businesses located in the U.S. are structured around the English language. Advertisements, menus, directions, etc. are often provided in English. Many of a business' employees may speak little or no foreign (that is, non-English) languages. Conversely, an increasing number of an employer's employees may have little or no skill in speaking English.\nIn order to provide effective service to people who do not speak English, a business currently has a limited number of options. First, it may recruit and hire employees who speak multiple languages and who are capable of communicating and serving speakers of languages other than English. This is problematic in that the number of available people in the job pool are correspondingly decreased and the cost of employing the employees is increased.\nSecond, a business may train its existing employees in other languages to enable them to better serve customers who speak languages other than English. This option is problematic in that current foreign language systems can be very costly and very time consuming, and aren't focused toward teaching the basic language skills needed to communicate on a basic level with a speaker of the second language. Also, most conventional training systems are focused primarily on teaching second language skills which encompass a large array of situations in which a person may require or use the second language, rather than focusing on specific needs for a particular profession.\nThese problems also arise in the instance where a business employs employees who speak little or no English, but who are nonetheless expected to serve English speaking customers."}
{"text": "Nowadays, in order to improve spectral efficiency, modern wireless communication systems use modern modulation and coding scheme techniques that involve multiple subcarriers direct or indirect mapping, such as OFDM (orthogonal frequency division multiplexing), OFDMA (orthogonal frequency division multiple access), SC-FDMA (single-carrier frequency division multiple access), WCDMA (wideband code division multiple access) or multi-carrier GSM (global system for mobile communication). These systems however suffer from high peak to average power ratios (PAPR), which originate from superposing subcarriers or various signals from diverse users. It is known that such kind of high PAPR is undesired because it requires more back-off for the subsequent power amplifier, and thus has a direct consequence on the amplifier's average output power. In order to reduce the dynamics of the incoming composite antenna signal, techniques like Crest Factor Reduction (CFR) are used to improve the PAPR and improve the efficiency of the transmission chain. However such processing algorithms may cause out-of-band spectral regrowth if not designed correctly, leading to non-compliance to regulatory spectral emission masks (SEM).\nThe noise shaping process is a clipping and filtering crest factor reduction technique which has been proposed to reduce the PAPR. The objective of the noise shaping process is to remove energy from the peaks of a signal and then to add that energy as a noise, inside and outside the signal bandwidth. As a result, the foregoing has the effect of reducing the PAPR as well as introducing in-band distortions which increase the bit-error ratio and also out-of-band degradations which interfere with communications in neighboring frequency bands.\nAlthough effective, it would be desirable to improve the noise shaping process to reduce the PAPR even more."}
{"text": "This invention relates to sensors that measure torque applied to a shaft. It particularly relates to sensors having a sleeve torsionally engaged with the shaft and sensing means responsive to torsional strain in the sleeve.\nMany known torque sensors operate by responding to magnetostrictive effects resulting from strain in a stressed member or transducer. Some of these are in commercial production. Efforts have been directed toward using magnetostrictive effects to measure the torque applied to the steering wheel by the driver of a motor vehicle. One known design torsionally engages a sleeve having desirable magnetostrictive properties to a portion of the steering wheel shaft. Another design uses the magnetostrictive properties of a current production steering wheel shaft to eliminate the cost of attaching a sleeve to the steering wheel shaft. In a third known design magnetostrictive material is beam or vapor deposited on a steering wheel shaft. The known designs have not proved entirely satisfactory. Known methods of attaching a sleeve require processes that are not easily adapted to large volume production. The same is true of beam or vapor deposition of magnetostrictive materials. Efforts to use the shaft itself have suffered from the difficulty of obtaining shafts consistently having desired magnetostrictive properties. For measuring steering torque in an automobile the ideal sensor would be inexpensive and compatible with existing steering wheel shafts.\nThe expression xe2x80x9ctorsionally engagedxe2x80x9d is used herein to describe engagement between a first element and a second element for transmitting torque therebetween. It includes engagement for transmitting torque by a rigid attachment such as a weld or adhesive joint or both elements being made of one piece of material. It also includes engagement by means that transmit only torque exemplified by a wrench socket engaging the head of a bolt. The expression xe2x80x9ctorsionally engagedxe2x80x9d is used to cover a broad range of torque transmitting engagement means that may or may not transmit forces in addition to torque.\nA torque sensor incorporating a sleeve of magnetostrictive material is described in U.S. Pat. No. 5,351,555 issued Oct. 4, 1994 to Garshelis. Particular attention is focused on the Garshelis patent because it is believed to offer the lowest cost sensor responsive to torque applied to a magnetostrictive sleeve. However, the invention is applicable to any torque sensor having a sleevelike transducer that is torsionally stressed when torque is applied to a shaft.\nThe Garshelis design provides a sleeve (xe2x80x9ctransducerxe2x80x9d) permanently magnetized in its circumferential direction. Garshelis discusses attachment of the transducer to the torsionally stressed shaft and (column 15 beginning at line 7) describes requirements which must be met by the chosen method of attachment:\nxe2x80x9cproper operation . . . requires that there be no slippage between any of the components at their interfaces. . . . Somewhat less obvious, but no less important, is the requirement that there be no inelastic strain in shaft 8 in any cross section which includes the transducer 4. Thus, all strains associated with the transmission of torque must be fully recoverable when the torque is relaxed.xe2x80x9d\nand in column 16 beginning at line 5\nxe2x80x9cAs already indicated, the transducer 4 and underlying shaft must act as a mechanical unit. Rigid attachment of the transducer 4 either directly or indirectly to shaft 8 is crucial to proper operationxe2x80x9d.\nIn fact, attachment by adhesive bonding (using known adhesives and known designs) or interference fit (Garshelis\"\" preferred method) do not satisfy the above quoted requirements. All known designs based on adhesive bonding or interference result in peak stresses exceeding the capabilities of the bond.\nIn column 16 beginning at line 5 and continuing through line 23 of column 17 Garshelis discusses three categories of torsional engagements between the transducer and the shaft. The categories are 1) salient point, i.e. splines, knurls, teeth etc. at the ends of the transducer mating with similar features on the shaft; 2) distributed, i.e. adhesive bonding or interference fit; 3) diffuse, i.e. welding or brazing the ends of the transducer to the shaft. The first xe2x80x9c1) salient pointxe2x80x9d and the last xe2x80x9c3) diffusexe2x80x9d work well but manufacturing methods for achieving these attachments are not easily adapted to automotive manufacturing procedures.\nAbout friction or adhesive bonding Garshelis states (column 16 lines 37 through 41):\nxe2x80x9cThis bonding limits the maximum measurable torque to a lower value than might otherwise be handled by the shaft 8 alone or transducer 4 alone, but is advantageous for other reasons as indicated previously.xe2x80x9d\nAccordingly, Garshelis expresses a known need for an xe2x80x9cadvantageousxe2x80x9d process such as adhesive bonding that does not limit the maximum measurable torque to xe2x80x9ca lower value than might otherwise be handled by the shaft 8 alone or transducer 4 alonexe2x80x9d. Garshelis goes on to state (column 16 lines 41 through 47):\nxe2x80x9cPress or shrink fits can be used to obtain the desired circular anisotropy, and can provide very substantial gripping forces which as a practical matter will not be broken by expected torques on shaft 8. With clean, degassed (and perhaps deoxidized) surfaces, the effective coefficient of friction can rise without limit and act somewhat like a weld.xe2x80x9d\nProviding xe2x80x9cclean, degassed (and perhaps deoxidized) surfacesxe2x80x9d on the elements before they are joined by press or shrink fits is expensive and time consuming. It is difficult to assure such qualities in many millions of parts as required for automotive production. It is not stated in the Garshelis patent but it is believed that to achieve in a press fit an effective coefficient of friction that xe2x80x9ccan rise without limit and act somewhat like a weldxe2x80x9d as stated in Garshelis the xe2x80x9cclean, degassed (and perhaps deoxidized) surfacesxe2x80x9d must be joined and heat treated at high temperatures in a suitable atmosphere for many hours. To obtain a shrink fit heat treatment is believed to be required both to achieve an effective coefficient of friction that xe2x80x9ccan rise without limit and act somewhat like a weldxe2x80x9d and to cause the shrinkage required for a shrink fit.\nAnother method of achieving an interference fit between the transducer and the shaft is described by Garshelis with reference to FIGS. 14, 15 and 16. In this method the shaft is hollow and an expander is drawn through the shaft to expand it thereby providing the desired hoop stress. This process also is believed to be difficult and expensive to implement in mass production of steering wheel shafts.\nThe following numerical examples will clarify the issues related to attaching a sleeve by adhesive or interference fit (without heat treatment or other processes to achieve an effective coefficient of friction that xe2x80x9ccan rise without limit and act somewhat like a weldxe2x80x9d). In column 10 lines 3 through 5 Garshelis cites the example of a shaft diameter of 0.5 inch (1.27 centimeters) and a transducer wall thickness in the 0.030 to 0.050 inch (0.076 centimeters to 0.127 centimeters) range. The wall thickness is important to achieve sufficient magnetic flux (Garshelis column 10 lines 24 through 31). From the well known fact that torque transmitted by a shaft is distributed as the third power of the radius it follows for the case of the aforementioned 0.5 inch diameter shaft that if the transducer and shaft have similar shear moduli (which is likely to be the case) 36 percent of the total torque will be transferred to the transducer in the case of 0.030 inch transducer wall thickness and 52 percent of the total torque will be transferred to the transducer in the case of 0.050 inch transducer wall thickness. A possible diameter of a steering wheel shaft of an automobile is 2 cm and it might be subjected to a maximum torque of 600 newton-meters (450 ft-lbs). Such a torque might be applied by a large healthy male driver after the wheel reached the end of its travel. At the one centimeter radius of the outer surface of the steering wheel shaft 600 newton-meters torque creates a tangential force of 60,000 newtons (13500 lbf). If 36 percent of the torque is transmitted to the transducer 21,600 newtons (4860 lbf) must be transferred between the transducer and the shaft by the attachment means. The fraction of the force transferred between the transducer and the shaft would be 36 percent in the case of the 2 centimeter diameter shaft if the inside diameter of the transducer is also 2 centimeters and the thickness of its wall is 1.2 millimeters (0.047 inches). The fraction would be much larger if the inside diameter of the transducer is larger and the thickness remains 1.2 millimeters. Assuming an adhesive shear strength of 10 newtons per square millimeter (1419 psi) and assuming means exist for providing constant shear stress over the area of adhesive attachment, transferring 21,600 newtons requires 21.6 square centimeters or 3.5 centimeters of shaft length of bonded area at each end of the transducer.\nThe second example is an interference fit. If the transducer wall thickness is 1.2 millimeters and is stressed to a hoop stress of 700 mpa (100,000 lbf/in2) and the coefficient of friction is 0.3, 1575 newtons (353 lbf) of shear force can be transferred per millimeter of length. Transferring the aforementioned 21,600 newtons of shear force requires about 1.4 centimeters of shaft length of contact with the shaft at each end of the transducer.\nIn summary, in the case of a two centimeter diameter steering wheel shaft, both bonding by adhesive and attachment by press fit would require contact with the shaft for one to four axial centimeters beyond each end of the active area of the transducer to transmit the forces encountered in operation assuming uniform shear forces. To prevent higher stresses that would cause adherence to fail the shear force must be distributed uniformly over the area of attachment. In fact, known technology does not enable the hereinabove reproduced requirements (Garshefis column 16 lines 37 through 41 and column 16 lines 41 through 47) to be achieved with any amount of adhesive or conventional press fit adherence area because the shear stresses peak at the ends of the attachment regions and exceed the maximum shear capabilities of adhesives and/or press fits.\nA substantial difference is now evident between xe2x80x9cdistributed attachmentxe2x80x9d (adhesive, friction) and xe2x80x9csalient point attachmentxe2x80x9d and xe2x80x9cdiffuse attachmentxe2x80x9d. In the latter two attachment is truly at the ends of the transducer and the transducer operates as a unit with the shaft. This is also true in the aforementioned case where the effective coefficient of friction rises without limit and acts somewhat like a weld which is believed to be properly categorized as a xe2x80x9cdiffuse attachmentxe2x80x9d. In the xe2x80x9cdistributed attachmentxe2x80x9d cases attachment forces are required to be distributed over lengths of shaft such as the aforementioned one to four centimeter attachment regions at each end of the transducer.\nIt will also be appreciated from the above numerical examples taken with the following that where the transducer has a constant thickness as illustrated in FIGS. 1, 3, 4 and 6 through 16 of Garshelis (all of the figures that illustrate transducers) the end portions of the transducer do not xe2x80x9cact as a mechanical unitxe2x80x9d with the steering wheel shaft unless the ends are effectively welded to the shaft. For the transducer to xe2x80x9cact as a mechanical unitxe2x80x9d with the steering wheel shaft it must twist as the steering wheel shaft twists over its entire length. However, if a constant thickness transducer is attached by adhesive or press fit, sufficient torque to twist the transducer and cause it to xe2x80x9cact as a mechanical unitxe2x80x9d with the steering wheel shaft is only achievable in a xe2x80x9ccentral regionxe2x80x9d between the aforementioned attachment regions. Outside the xe2x80x9ccentral regionxe2x80x9d the torque available to twist the transducer diminishes with distance from the xe2x80x9ccentral regionxe2x80x9d because the transmission of torque is xe2x80x9cdistributedxe2x80x9d and the twisting of the transducer diminishes as the torque diminishes with distance from the central regions causing the torsional strain of the transducer and the shaft to be different far from the central regions. In Garshelis\"\" words cited hereinabove: xe2x80x9cThis bonding limits the maximum measurable torque to a lower value than might otherwise be handled by the shaft 8 alone or transducer 4 alone.xe2x80x9d\nAn object of this invention is to provide a torque sensor transducer which can be attached by adhesive to a torque carrying shaft and which will then operate xe2x80x9cas onexe2x80x9d with the torque carrying shaft.\nA general object of this invention is to provide a torque sensor which also overcomes certain disadvantages of the prior art.\nThe present invention provides a torque sensor for measuring the torque applied to a shaft. It comprises a magnetostrictive sleeve torsionally engaged with two shear levelers. The shear levelers are bonded by adhesive to the shaft. The shear levelers have flared ends and regions of varying torsional elasticity that operate to level the shear stress in the adhesive. The term xe2x80x9clevelxe2x80x9d is used herein with reference to shear stress in adhesives to describe causing the shear stress to be constant and without peaks over the area bonded by adhesive. It may include being constant at two or more different levels at two or more areas bonded by adhesive.\nFurther, in accordance with the invention, the torque sensor is attached to the shaft by adhesive which is stressed in shear without stress peaks thereby enabling designs wherein the adhesive can transfer torques approaching the yield limit of the shaft.\nFurther, in accordance with the invention, the shear levelers have varying torsional stiffnesses to provide a uniform shear stress in the adhesive.\nFurther, in accordance with the invention, the shear levelers have flared ends and varying thickness adhesive at the flared ends further distributes stress in the adhesive and enables designs wherein the adhesive transmits torques that approach the yield torque of the shaft.\nFurther, in accordance with a first embodiment of the invention, low magnetic permeability isolation rings magnetically isolate the shear levelers from the magnetostrictive central segment.\nFurther, in accordance with the aforementioned first embodiment of the invention, the isolation rings are welded to the shear levelers and the magnetostrictive central segment.\nFurther, in accordance with the aforementioned first embodiment of the invention, the magnetostrictive central segment is pressed onto a stack of washers with crowned outer circumferences that maintains the transducer in its cylindrical shape and minimizes the torque that must be accumulated by the shear levelers. Great hoop stress in the magnetostrictive central segment is achieved by heat treatment after the magnetostrictive central segment is pressed onto the stack of washers.\nFurther, in accordance with the aforementioned first embodiment of the invention, each washer of the aforementioned stack of washers with crowned outer circumferences is coated with a thin layer of material that evaporates during heat treatment thereby leaving each washer separated from adjacent washers and therefore free to rotate without friction when the transducer is torsionally strained.\nFurther, in accordance with a second embodiment of the invention, the shear levelers are unitary with a low magnetic permeability middle segment upon which the magnetostrictive central segment is pressed and welded and possibly shrunk whereby great hoop stress in the magnetostrictive central segment is achieved which advantageously provides desirable magnetic properties.\nFurther, in accordance with the aforementioned second embodiment of the invention, the shear levelers are unitary with a low magnetic permeability middle segment upon which the magnetostrictive central segment is placed and welded and great hoop stress in the magnetostrictive central segment is achieved by expanding the middle segment and the magnetostrictive central segment together which advantageously provides desirable magnetic properties.\nFurther, in accordance with a third embodiment of the invention, the shear levelers are unitary with the magnetostrictive central segment and annular grooves are provided between the shear levelers and the magnetostrictive central segment. The grooves enhance magnetic anisotropy and provide surfaces against which force may be applied to facilitate installation of the transducer on the shaft.\nFurther, in accordance with the invention, a torque sensor comprises a circularly symmetric magnetic element centered on the rotation axis of a magnetostrictive element for providing a lower reluctance magnetic field path and less sensitivity to a bent shaft or other asymmetry.\nA complete understanding of this invention may be obtained from the description that follows taken with the accompanying drawings."}
{"text": "Existing communications methods and systems are overly power hungry and/or spectrally inefficient. Further limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such approaches with some aspects of the present method and system set forth in the remainder of this disclosure with reference to the drawings."}
{"text": "Currently, the Wi-Fi Alliance is proposing that packets be aggregated into larger groups of packets before being sent in tunnel mode to a remote station (STA). For stations that support a VHT (Very High Throughput) mode, the number of aggregated packets may be very large. The aggregation is a key part of reaching a stated goal of 1 gigabit per second for wireless communications.\nThe selection of the number of packets to be aggregated greatly affects the transmission rate, the amount of overhead, and the airtime required to transmit the packet. In addition, to aggregate packets, the system must wait for all of the packets that are to be aggregated to arrive before they can be aggregated. In addition, the system must wait for enough packets to arrive to satisfy the conditions for aggregation. This introduces delays that impact applications or user experiences."}
{"text": "Certain oxadiazine compounds have been described as useful as pesticides and as pharmaceutical agents. For example, U.S. Pat. No. 5,536,720 describes substituted 2-phenyl-1,3,4-oxadiazine-4-carbamide compounds useful as insecticides and acaricides. Trepanier et al, J. Med. Chem 9: 753-758 (1966) describe certain 2-substituted 4H-1,3,4-oxadiazines useful as anticonvulsants in mice. U.S. Pat. No. 3,420,826 describes certain 2,4,6-substituted 4H-1,3,4-oxadiazines, useful as sedatives, anticonvulsants, and as pesticides against nematodes, plants, and fungi. U.S. Pat. No. 3,420,825 describes methods for producing certain 2,4,6-substituted 4H-1,3,4-oxadiazines.\nIt is a purpose of this invention to provide novel dihydrooxadiazine derivatives useful as insecticides."}
{"text": "This section is intended to provide information relevant to understanding various technologies described herein. As the section's title implies, this is a discussion of related art that should in no way imply that it is prior art. Generally, related art may or may not be considered prior art. It should therefore be understood that any statement in this section should be read in this light, and not as any admission of prior art.\nIn modern manufacturing of integrated circuits, conventional cell architecture designs are typically generated with two-dimensional Electronic Design Automation (EDA) tools that are inefficient and cumbersome to employ. Some 2D design flows may directly adopt a 2D placement solution to determine the standard cell and inter-tier via locations, with limited results for three-dimensional (3D) optimizations. Some other 2D design flows propose block partitioning for 3D optimizations across multiple tiers. However, the inter-tier via locations are typically determined by the 2D placement solution from a single tier, which lacks co-optimization between multiple tiers."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus which displays a preview upon outputting images of a plurality of pages, a control method thereof, and a storage medium.\n2. Description of the Related Art\nSome image forming apparatuses such as digital multifunction peripherals (MFPs) include a function of displaying a thumbnail of a generated page image as an image formation target on a liquid crystal panel of an operation unit as a preview so as to prevent a wasteful output due to a setting error. In an image forming apparatus which includes post-processing functions such as stapling, folding, and punch hole functions, a method of processing a page image generated based on input print data, and displaying a finish preview which allows the user to recognize a final output form applied with post-processes has been proposed.\nFor example, Japanese Patent Laid-Open No. 2003-87560 has proposed an image forming apparatus which displays pages in a two-page spread state except for the first and last pages, when post-processing is stapling processing upon displaying a finish image indicating a state applied with the post-processing by processing image data obtained by converting input image data to reduced-scale data. Also, Japanese Patent Laid-Open No. 2010-35030 has proposed a preview method which prevents user's confusion by displaying finish information outside a print sheet area of a preview image so as to improve visibility, and deciding a direction of the preview image in accordance with the convey direction of a sheet.\nHowever, the above related arts suffer the following problems. For example, in an image forming apparatus which includes a preview function in consideration of post-processing like in the related art, the user can confirm if post-processing settings are as the user intended, but a configuration which allows the user to recognize a relationship between pages before and after a preview target page is not provided. That is, in a preview of the related art, a preview image is displayed by reflecting a finish result for each page image of a sheet, or only the method of displaying sheets in a two-page spread state upon designation of stapling as a case in which pages are displayed to allow easy understanding of an anteroposterior relationship has been proposed. For this reason, in an image forming apparatus including various finishing functions and output instruction forms, a preview method which allows the user to easily recognize a relevance according to output forms between a plurality of pages such as the anteroposterior relationship or positional relationship of a plurality of pages to be output is demanded."}
{"text": "This invention relates to a method for the cleavage of carboranyl carbon-silicon bonds in compounds and polymers containing same, and in addition, to a method for preparing meta-carborane using this technique.\nORTHO-Carborane, of the empirical formula C.sub.2 B.sub.10 H.sub.12, has an icosahedral structure and was the first icosahedral carborane discovered and prepared in quantity. Subsequently, methods were developed for isomerizing the ortho derivative to the meta-and para-isomers. For example, ortho-carborane may be rearranged or isomerized to meta-carborane in a hot tube flow process at high temperatures (i.e., 600.degree. C) involving short residence times. See S. Papetti, et al, Industrial Engineering Chemical Product Research and Development, 5,334 (1966). At these temperatures, carborane vapors may burn or explode if allowed to contact air and hence such a process is not attractive for industrial operation.\nAnother process is known for isomerizing carboranes in bulk and is disclosed in U.S. Pat. No. 3,440,265. The process involves the thermal rearrangement of large silyl group-substituted ortho-carboranes to the metal-isomer at significantly lower temperatures than the prior technique. In contrast to the earlier flow process, this bulk isomerization occurs readily at the reflux temperature and can be performed using conventional laboratory equipment. Furthermore, the method is adaptable to large-scale industrial production.\nThe metal-isomer has become increasingly important as the key ingredient of high performance metal-carborane-siloxane polymers. Such polymers have excellent high temperature characteristics and can be prepared by methods known in the art. For example, linear, high molecular weight carborane-siloxane polymers and methods for their preparation are disclosed in commonly-assigned, copending U.S. Pat. Application Ser. No. 770,509, filed Feb. 22, 1977. These polymers may be compounded with suitable fillers and additives and fabricated into cured elastomeric articles. The excellent high temperature properties of such elastomers make them especially useful for a wide variety of applications, such as gaskets, seals, wire and cable insulation, and the like. Although carborane-siloxane polymers have excellent high temperature and other high performance characteristics, they are extremely expensive to produce. Hence, it would be an advantage to be able to reclaim the carborane portion from scrap, off-grade or otherwise used-up polymer, and also to recover or regenerate carborane monomer by-products used in forming the carborane-siloxane polymers.\nIt is known in the art that the carborane icosahedral structure is subject to attack by strong bases. See, for example, Papetti, et al, Inorganic Chemistry, 3 1444 (1964) which discloses without detailed explanation that some cyclic silylated carboranes are cleaved. In addition, Schwartz, et al, Inorganic Chemistry, 4, 661-4 (1955) indicate that a very labile diethoxysilylated carborane is cleaved under acidic or basic conditions not specified.\nThe present invention provides a technique for highly selectively cleaving carboranyl carbon-silicon bonds in monomeric silylated carboranes and carborane polymers. The method taught herein is highly selective in that the carborane cage contained in these materials may be recovered intact which is highly advantageous due to the expensive nature of the materials involved. Further, this novel technique may be employed in a process for the formation of metal-carborane in bulk when used in combination with known techniques. This represents a significant advance since this method for preparing meta-carborane in high yield in bulk may be practiced on a large-scale, industrial level."}
{"text": "This invention relates to an air conditioning system and more particularly to those having freeze tanks using ice to store cooling power.\nGenerally, in the past, devices have been used wherein a coolant is partially frozen and the part remaining liquid is used as a coolant. An article describing one of these systems is shown in the September, 1979 issue of Mechanix Illustrated.\nPatents which are related to devices of this type are set forth below: U.S. Pat. No. 2,097,556; U.S. Pat. No. 2,158,707; U.S. Pat. No. 2,448,453; and U.S. Pat. No. 2,737,027."}
{"text": "It is a problem with currently available solutions that binaural transmission, i.e. transmission between the two hearing instruments of a binaural hearing aid system, creates an additional latency to the hearing aid processing because of signal buffering, quantization, coding, synchronization, etc. Typically, the hearing instruments may either only transmit or receive audio signals at a given time instant. This means that in order to have hearing instruments operate synchronously within a time slot, the hearing instruments are required to wait until the audio signal packages have been transmitted and received at both hearing instruments.\nBy only transmitting the audio signal in one direction, such processing delay may be reduced, because waiting for an audio package, which is transmitted in the opposite direction will not be necessary. Therefore, in order to design binaural signal processing algorithms (e.g., binaural noise reduction algorithms) which make use of signals sent only in one direction, there is a need to provide a solution to the problem of how to decide, at any given moment, in which direction (i.e. from which hearing instrument to which hearing instrument) the signal is to be sent."}
{"text": "Universal Mobile Telecommunications System (UMTS) networks have seen an explosive data growth in past few years and, in future, are expected to see continuing growth in the Packet Switched (PS) domain. Beyond data traffic volume growth, an even more aggressive growth in data signaling load has been detected. Among all the signaling messages/procedures on UMTS networks, Radio Access Network (RAN) signaling procedures have caused the most growth and impact. This is due to complicated radio resource sharing techniques required to conserve resources occupied by various users and services.\nThe majority of RAN signaling events are for connection setup and state transitions (e.g., during Channel Switching). Typically, when a data payload is to be sent from/received by a user equipment (UE), a request is sent to a radio network controller (RNC) to establish a dedicated channel (DCH). Once the data payload is sent or received, multiple inactivity timers are triggered by the RNC and upon expiration of the timers, the RNC transitions the UE from DCH to forward access channel (FACH) and then to IDLE state. To achieve resource efficiency, such timers are often set to short values (cumulatively around 12-16 seconds). Thus, the UE is quickly moved into the IDLE state after completion of a current data session (download and/or upload). Since there is no active data connection between the UE and the core network during the IDLE state, power consumption is minimized. Data applications on the UE initiate multiple data payloads for communication between the UE and RNC. However, the applications operate independently from the radio network perspective, which leads to requesting and establishing multiple independent connections for payloads from different applications. As a result, even though battery life of the UE is conserved, a large number of signaling events are generated and RNC processing load is substantially increased.\nIn addition to network initiated inactivity-based state transition, UE manufacturers have introduced a fast dormancy (FD) feature that initiates direct transition from DCH to IDLE or FACH to IDLE, before the respective network inactivity timer expires. In this type of system control, the UE proactively releases the data connection, established by the RNC, directly from DCH to IDLE or FACH to IDLE as quickly as possible, to further conserve UE battery life. However, once the UE is in the IDLE state, the data connection must be reestablished to communicate another payload. The reestablishment of the data connection is resource intensive, consumes a high amount of power in the RNC, and can significantly drive up the RNC load."}
{"text": "Particularly for medical purposes, an image which is stored in a storage layer as a latent image is produced from an object, for example a patient, by means of X-ray radiation. It is mostly a phosphor layer which is used as the storage layer in this case. The storage layer is excited by means of a radiation source in order to read out the X-ray image stored in the storage layer. Therefore, on the basis of this excitation the storage layer emits light which has an intensity corresponding to the stored X-ray image. The light emitted by the storage layer is received by a receiving means and subsequently converted into electric signals such that the X-ray image stored in the storage layer can subsequently be visualized. The X-ray image can be displayed, for example, directly on a monitor, or else be written onto a photographic X-ray film which can be used specifically for X-ray images. After the X-ray image has been read out from the storage layer, the latter is erased by an erasing device, in order to be able to store a subsequent X-ray image.\nSuch storage layers can be arranged in an X-ray cassette. European patent EP 0 288 014 B1 discloses erasing a storage layer which is located in an X-ray cassette. The X-ray cassette contains an opening mechanism such that the X-ray cassette can be opened in order to erase information stored in the storage layer. Present outside the X-ray cassette is an erasing means which outputs an erasure radiation. This erasure radiation is irradiated into the open X-ray cassette and projected onto the storage layer. The irradiation of the storage layer by means of the erasure radiation erases the information stored in the storage layer. The known X-ray cassette contains a reflector or a diffusing means for the purpose of reflecting the erasure radiation irradiated into the cassette, and of scattering it over the width of the storage layer.\nA device for reading out information stored in a storage layer is disclosed in patent U.S. Pat. No. 5,038,037. An X-ray table for X-raying patients is described therein. The X-ray table contains two phosphor layers which are fitted on the top side and underside of a tape and serve for storing X-ray information. Both for the phosphor layer fitted on the top side and for that fitted on the underside of the tape, the known X-ray table has a dedicated device for reading out information from this phosphor layer. Furthermore, the X-ray table contains two erasing means which are used to erase the phosphor layers. A dedicated erasing means is provided for each of the phosphor layers. Serving as erasing means is an erasing lamp which—viewed in the transport direction of the tape—is arranged upstream of that point at which the X-ray image of the patient is taken. As a result, an “old” X-ray image stored in the phosphor layer is erased before a “new” X-ray image is taken. A plurality of erasing lamps arranged next to one another are likewise used as erasing means, being arranged in the X-ray table below that position at which the X-ray image is projected onto the stationary tape with the phosphor layers. Before the taping of a “new” X-ray image, the “old” X-ray image is therefore firstly erased with the aid of this erasing means, the tape being stationary. In this type of erasing means, the phosphor layer cannot be used directly for a “new” picture. The erasure of the phosphor layers does not take place until immediately before the subsequent taking of a “new” X-ray image. This means that time is lost and tuning takes place between the operations of erasure and subsequent storage. The erasure must firstly be terminated before the next X-ray image can begin to be taken.\nPatent application WO 99/28765 discloses an X-ray cassette in which both a device for reading out information stored in the storage layer and an erasing means for erasing information stored in the storage layer are present. The erasing lamp extends over the entire width of the storage layer in which information can be stored. By means of a drive, the erasing means is guided over the storage layer along a transport direction running perpendicular to the line direction. The storage layer can be erased in this way line by line."}
{"text": "In paper manufacturing processes, cellulose fibers are deposited as an aqueous slurry on a screen. The water is removed to form the paper. Frequently, a dewatering fabric is used in a second stage of water removal from the aqueous slurry.\nEschmann (U.S. Pat. No. 5,204,171) describes a papermaking fabric comprising a blocking layer.\nChuang et al (WO9616305) describes a capillary dewatering method and apparatus.\nNew fabrics for dewatering during the papermaking process are desired because improved efficiency in water removal can provide cost savings by reducing or eliminating the need for subsequent elevated temperature drying steps. The present invention is directed to these and other important ends."}
{"text": "With the development of multimedia technologies, in an increasing number of application scenarios, multiple devices need to be used to constitute a synchronous playback system for synchronous audio playback, to implement synchronization of sounds played from the multiple devices, so as to implement a function and an effect that are difficult to be implemented by a single device. For example, synchronous audio playback performed by a synchronous playback system constituted by multiple mobile phones can implement effects that cannot be implemented by a single mobile phone, such as multi-channel surround sound and mobile phone speaker array.\nIn the prior art, an audio playback manner generally used by a synchronous playback system is controlling devices in the synchronous playback system to play an audio file simultaneously. The prior art has the following characteristic, when an audio file is played in a synchronous playback system, the audio file is generally divided into several audio clips, and then the audio clips are played one by one. However, playback of each audio clip requires processing procedures such as decoding, audio mixing, digital signal processing (DSP), and output, and therefore, there is a time difference from start of playback of an audio clip to actual generation of a sound, and this time difference is a playback delay. A value of the playback delay is determined by a hardware configuration and a software environment of a device, and when hardware configurations or software environments of various devices are different, values of playback delays are also different.\nBecause playback delays of various devices are different, when various devices start to play a same audio clip simultaneously, sounds are actually generated by the various devices at different time points. When hardware configurations or software environments between various devices differ greatly, a time deviation in actual generation of sounds also becomes large. For example, in a case in which a HUAWEI MATE1 mobile phone and a HUAWEI HONOR mobile phone start to decode a same audio clip simultaneously, a time difference in actual generation of sounds may be 130 milliseconds (ms) to 168 ms. Further, a time deviation between various devices is generated when each audio clip is played, and with an increasing number of played audio clips, the time deviation in actual generation of sounds between various devices gradually accumulates. As a result, the problem of time inconsistency between various devices in actual generation of sounds becomes increasingly serious, which makes it difficult to implement synchronous audio playback."}
{"text": "The present invention is directed toward a portable golf practicing device and more particularly toward a golf club having a telescoping shaft and a ball retrieval system such as a fishing reel attached thereto.\nGolf clubs with telescoping shafts have been proposed in the past. One such golf club is shown, for example, in U.S. Pat. No. 3,424,464. This patent shows a golf club in the form of a putter which has a telescoping shaft and which is intended primarily as a novelty item. The putter is mounted on a small stand which has an open front and can, therefore, function as a practice cup. Thus, the putter can actually be used by an executive or the like practicing putting in his office.\nA golf club with a telescoping shaft is also shown in U.S. Pat. No. 2,107,983. This patent is directed toward a club which is intended to be actually used on a golf course and which can be adjusted during the course of play to the effective length desired without materially disturbing the balance of the club. This patent, as with the prior patent described above, is directed toward a golf club only and includes no means for retrieving a ball after it has been hit.\nGolf clubs have also been proposed which include retrieval devices in combination therewith so that a ball can be retrieved after it has been hit. These prior devices include fishing reels attached near the handle of the golf club as shown, for example, in U.S. Pat. Nos. 3,065,563 and 4,526,374.\nThe first-mentioned patent includes a fishing reel which is relatively permanently fixed to the golf club handle and the shaft of the golf club includes a plurality of eyelets secured thereto for guiding the line. The second patent, on the other hand, includes a fishing reel carried by its own bracket which is adapted to be secured to substantially any golf club. In both cases, the end of the line is secured to a golf ball.\nWhile each of the devices shown in U.S. Pat. Nos. 3,065,563 and 4,526,374 show the use of a retrieval device in the form of a fishing reel, the devices are used on substantially conventional golf clubs. There is no suggestion in these patents that the golf club be provided with a telescoping shaft. Furthermore, the fishing reels must be manually rewound in order to retrieve the ball after it has been hit. No automated means for retrieving the ball is proposed."}
{"text": "In the art of deposition of films of material onto a surface of a substrate, there are many known techniques, including vacuum evaporation deposition, ion plating, ion- and plasma-assisted sputtering or chemical vapor deposition (\"CVD\"), and the more modern ICB approach. ICB deposition is an ion-assisted technique in which the material to be deposited on a substrate is heated in a crucible and its vapor ejected through a small nozzle into a vacuum region. The vapor forms loosely-held atomic clusters, each cluster comprising 100 to 2000 atoms of the material. Some of the ejected vaporized atomic material is ionized by electron bombardment and the atoms are accelerated toward the substrate disposed in the vacuum region. The ionized material, together with the neutral (i.e., non-ionized) component of the vapor, arrive at the substrate surface for deposition thereon. ICB deposition offers the ability to precisely control the deposited film structure by applying kinetic energy to the vapor clusters during film deposition. Kinetic energy control is achieved by varying the acceleration voltage and the electron current for ionization.\nIn most of the known ion- and plasma-assisted deposition techniques, the individual atoms of the material to be deposited on the substrate generally impact the substrate surface with excessive kinetic energy, producing a relatively high number of defects in the substrate and/or the deposited film. With ICB deposition, a more useful lateral energy is obtained as the clusters impact the substrate and the atoms break off, without damaging the film and substrate. Due to the effects of ionized cluster bombardment, ICB deposition produces films with high density, strong adhesion, a low impurity level, and a smooth surface. Also, film properties usually associated with relatively high substrate temperatures in conventional vacuum depositions can be obtained at lower substrate temperatures in the ICB technique. This results in a distinct advantage in semiconductor device fabrication. See, for example, U.S. Pat. Nos. 4,152,478, 4,217,855, 5,350,607 and 5,380,683, all of which are hereby incorporated by reference.\nThe trend in the semiconductor manufacturing industry has always been to increase the number of active devices (e.g., transistors, resistors, capacitors) formed in an area of a semiconductor substrate (e.g., silicon, germanium, gallium arsenide). This increase in IC density has been achieved primarily by decreasing the size of the active devices and associated electrically-isolating areas (e.g., field oxide) formed within the IC substrate. Sizes of such devices and areas are now down into the submicron range. This size reduction has been achieved largely through improved fabrication methodologies and structures.\nFor example, it is common to employ one or more thin horizontal layers of a high-conductivity metal (e.g., aluminum), separated by insulating layers, together with associated vertical, high-conductivity interconnecting plugs or vias routed through the insulating layers. This is to connect between the active devices in the semiconductor substrate. In this approach, holes are formed in the insulating layers at desired connection points. The holes are filled with a high-conductivity, low-resistivity material to form the vias. The overlying horizontal metallization layer is then deposited. This process is repeated for subsequent vertical plugs and horizontal overlying layers. In this way, an ohmic contact is formed between active device areas (e.g., the source and drain terminals of an MOS transistor, or the collector, base and emitter terminals of a bipolar transistor) in the silicon substrate. Interconnections can be made in the different metal layers as cross-overs, thereby further increasing IC density. This interconnect process is generally referred to as \"metallization\".\nHowever, conventional metallization techniques encounter problems as device geometries become increasingly smaller. For example, as device sizes shrink, contact or via holes formed closer together, are of smaller diameter and have steeper vertical sidewalls. As a result, it has become difficult to accurately deposit conventional materials, such as aluminum, using conventional techniques (e.g., sputtering), into the holes to achieve uniform contact with device contact areas and metallization layers. Non-uniform contact within the interconnects results in problems such as non-planar topographies and electrical breaks formed at the edges of the holes. The result is manifested in relatively poor step coverage of the metallization layers and vias. Other problems include those related to the reliability of the interconnect lines, such as electromigration or wear-out failures.\nAlso, as device sizes shrink, connection lines become smaller, thereby subjecting the interconnects to higher current densities. At these higher densities, aluminum metallization layers and plugs are increasingly susceptible to stress migration and electromigration damage. Further, conventional metallization materials such as aluminum exhibit higher resistance, which increases the RC time constants of the device, thereby limiting overall device speed.\nIn an attempt to overcome these problems, it is known in the prior art to utilize tungsten as the material comprising the metallization layers and/or the plugs, or as a separate, additional barrier layer utilized in conjunction with an aluminum layer or plug. Tungsten is utilized in part because, relative to aluminum, tungsten has a higher resistance to electromigration. The tungsten barrier layer acts as a backup layer, maintaining the electrical connection integrity of the metallization layers and plugs if the primary aluminum layers and plugs fail due to stress migration and electromigration. Still further, tungsten can be deposited somewhat more uniformly than aluminum, thereby reducing the aforementioned concerns with step coverage.\nHowever, an inherent problem with the use of tungsten alone as the material comprising the metallization layers and plugs is that tungsten has approximately three times the resistivity of aluminum. Therefore, it is desired to utilize tungsten as the primary constituent of a tungsten alloy formed as a metallization layer or plug, and to deposit the tungsten alloy layer onto a semiconductor substrate using the ICB approach."}
{"text": "Securing items during transport from one location to another location can be difficult for those without a vehicle. Depending on the distance from the retailer, the act of transporting the items can be physically demanding and frozen items can melt or defrost along the way."}
{"text": "2-bromo-4,6-dinitromesitylene may be used as a starting material to form a number of compounds, such as for example 2-perfluoroalkyl-4,6-diaminomesitylene, which in turn may be used as one of the main constituents in polyamides, polyimides, urethanes, ureas, and the like.\n2-bromo-4,6-dinitromesitylene has been prepared so far by reacting nitric acid or a mixture of fuming sulfuric and nitric acids with bromomesitylene. Such preparations are shown for example in very old as well as more recent references, such as \"Fittig and Storer, Liebigs Ann. Chem., 147, 1-11 (1868)\", \"Suessenguth, Liebigs Ann. Chem., 215, 242-252 (1882)\", and \"Adams and Miller, J. Amer. Chem. Soc., 62, 53-56 (1940).\"\nThe following is translated from Fittig and Storer (p. 8): \"By pouring fuming nitric acid over bromomesitylene at ambient temperature and after short standing, it is converted into the (title) compound, which through washing with water and recrystallization with alcohol is easily purified. It crystallizes in fine, colorless needles, which dissolve little in cold alcohol, easier in hot alcohol, yet are much less soluble than the starting compound. Its melting point is 189.degree.-190.degree. C. Under warming, the compound develops a genuine musk smell.\"\n______________________________________ Analysis: Calculated Found ______________________________________ C9 108 37.37% 37.16% H9 9 3.11 3.29 Br 80 27.68 27. 2 NO.sub.2 92 1.84 -- 289 100.00 ______________________________________\nThis preparation appears to be simple and leading to a pure product, since there are only two positions in the aromatic ring of bromomesitylene, which are open for substitution by nitro groups. Simple elemental analysis, indicates the reaction to be clean and straightforward. However, more careful examination of this reaction reveals it to be considerably more complicated.\nDepending on how the reaction is conducted, a number of undesirable by-products, such as 2,4,6-trinitromesitylene, 2,4-dibromomesitylene, and 2,4-dibromo-6-nitromesitylene, are formed. Some of these by-products are extremely difficult to separate, because they cocrystallize with the desired product, even upon repeated recrystallization from a variety of solvents. These by-products may also prove to be catastrophic to polymers intended to be made from derivatives of the main product (dinitro compounds) after hydrogenation (diamines), since monoamine derivatives will serve as chain terminators, and triamine derivatives will serve as crosslinkers. Thus, expensive procedures, such as repeated crystallizations, sublimations, and the like are needed for purifying the products obtained in subsequent conversions of impure 2-bromo-4,6-dinitromesitylene, to free these products from undesirable impurities.\nIn contrast to previous attempts, considerably purer 2-bromo-4,6-dinitromesitylene may be prepared according to the methods of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to an air-controlled support for a vehicle seat and, more particularly, to an air-controlled lumbar support apparatus for use in a vehicle seat.\n2. Description of the Prior Art\nA conventional air-controlled lumbar support apparatus is disclosed in U.S. patent application Ser. No. 4,190,286 to Bentley. An example of the conventional apparatus is shown generally in FIG. 1 and includes a vehicle seat 10 with a backrest 12 having a series of inflatable chambers 14. Each of the inflatable chambers 14 is connected to a valve mechanism 16 by individual tubes 18. The valve mechanism 16 controls the intake and exhaust of the air within the inflatable chambers 14. Additional air is supplied to the inflatable chambers 14 by means of an air pump 20, which delivers the air supply through the valve mechanism 16.\nSince several inflatable chambers 14 are installed within the backrest 12 of the seat 10, the valve mechanism 16 requires a series of valves (not shown) for controlling the intake and exhaust of air to the inflatable chambers 14. As a result, the valve mechanism 16 typically is oversized and unduly complicated. Furthermore, the valve mechanism 16 and the air pump 20 must be independently mounted on the seat 10. Therefore, considerable space is required to position the tubes 18, the valve mechanism 16, the inflatable chambers 14 and the air pump 19 in reasonable proximity to each other."}
{"text": "1. Field of the Invention\nThis invention relates to display systems utilizing holographic lenses and more particularly to a display system incorporating a holographic lens constructed with aberrated wavefronts which is used in conjunction with other optical elements including cylindrical, tilted and decentered surfaces to provide correction of the aberrations in the holographic lens.\n2. Description of the Prior Art\nPreviously, optical systems utilizing holographic lenses have been designed with good image quality over small fields of view and poor image quality over large fields of view. Also, optical systems utilizing holographic lenses have been required to operate at relatively small off-axis angles because of the aberrations inherent in a holographic lens utilizing large off-axis angles. Aberrations arise for the on-axis image point in a holographic lens designed for a large field of view because the wavefront used in the reconstruction process is not identical to either construction beam. This arises because there is a pupil-to-pupil relationship between the construction beam sources while an object-to-image relationship must hold upon reconstruction. The holographic lens is constructed with coherent sources located close to or at the desired entrance and exit pupil locations of the holographic lens so that the Bragg condition is met for the chief rays, thus providing a high reconstruction efficiency across a large field of view. The principal aberrations in a holographic lens constructed as above and operating at a relatively large off-axis angle are axial astigmatism, axial coma and an assymetric variation of coma and astigmatism across the field. In addition there is a considerable amount of field tilt and distortion. The assymetric variation of astigmatism arises principally because the field tilt for fans of rays traced in a vertical plane is different from the field tilt for fans of rays traced in a horizontal plane. Unless compensation is provided for the aberrations in the holographic lens operating at a relatively large off-axis angle, the display provides a very poor image quality that may be unsatisfactory for operations such as helmet-mounted displays or aircraft head-up displays."}
{"text": "Magnetic discs and disc drives provide quick access to vast amounts of stored information. Both flexible and rigid discs are available. Data on the discs is stored in circular tracks and divided into segments within the tracks. Disc drives typically employ one or more discs rotated on a central axis. A magnetic head is positioned over the disc surface to either access or add to the stored information. The heads for disc drives are mounted on a movable arm that carries the head in very close proximity to the disc over the various tracks and segments.\nThe increasing demands for higher areal recording density impose increasingly greater demands on thin film magnetic recording media in terms of coercivity (Hc), remanent coercivity (Hr), magnetic remanance (Mr), which is the magnetic moment per unit volume of ferromagnetic material, coercivity squareness (S*), signal-to-medium noise ratio (SMNR), and thermal stability of the media. These parameters are very important to the recording performance and they depend primarily on the microstructure of the materials of the media. For example, as decreasing the grain size reduces the SMNR or reducing exchange coupling between grains, it has been observed that the thermal stability of the media decreases.\nThe requirements for high areal density, i.e., higher than 30 Gb/in2, impose increasingly greater requirements on magnetic recording media in terms of coercivity, remanent squareness, medium noise, track recording performance and thermal stability. It is extremely difficult to produce a magnetic recording medium satisfying such demanding requirements, particularly a high-density magnetic rigid disk medium for longitudinal and perpendicular recording.\nAs the storage density of magnetic recording disks has increased, the product of Mr and the magnetic layer thickness t has decreased and Hr of the magnetic layer has increased. This has led to a decrease in the ratio Mrt/Hr. To achieve a reduction in Mrt, the thickness t of the magnetic layer has been reduced, but only to a limit because the magnetization in the layer becomes susceptible to thermal decay and medium noise.\nMedium noise in thin films is a dominant factor restricting increased recording density of high-density magnetic hard disk drives, and is attributed primarily to inhomogeneous grain size and intergranular exchange coupling. Accordingly, in order to increase linear density, medium noise must be minimized by suitable microstructure control.\nLongitudinal magnetic recording media containing cobalt (Co) or Co-based alloy magnetic films with a chromium (Cr) or Cr alloy underlayer deposited on a non-magnetic substrate have become the industry standard. For thin film longitudinal magnetic recording media, the desired crystallized structure of the Co and Co alloys is hexagonal close packed (hcp) with uniaxial crystalline anisotropy and a magnetization easy direction along the c-axis that lies in the plane of the film. The better the in-plane c-axis crystallographic texture, the more suitable is the Co alloy thin film for use in longitudinal recording to achieve high remanance and coercive force. For very small grain sizes coercivity increases with increased grain size. The large grains, however, result in greater noise. Accordingly, there is a need to achieve high coercivities without the increase in noise associated with large grains. In order to achieve low noise magnetic recording media, the Co alloy thin film should have uniform small grains with grain boundaries capable of magnetically isolating neighboring grains thereby decreasing intergranular exchange coupling. This type of microstructural and crystallographic control is typically attempted by manipulating the deposition process, and proper use of underlayers and seedlayers.\nIt is recognized that the magnetic properties, such as Hcr, Mr, S and SMNR, which are critical to the performance of a magnetic alloy film, depend primarily upon the microstructure of the magnetic layer, which, in turn, is influenced by the underlying layers, such as the underlayer. It is also recognized that underlayers having a fine grain structure are highly desirable, particularly for growing fine grains of hcp Co alloys deposited thereon.\nFor high signal to noise ratio (SNR) magnetic recording media, it is desirable to have a high signal in a very thin film. Higher signal can be achieved by increasing the saturation magnetization (Ms) of the material at the top of the magnetic layer, and correspondingly increasing the fringing magnetic field that provides signal. Prior art magnetic recording systems generally employ media including a magnetic layer alloy including Co and Cr, and other elements often including Pt, and B. These magnetic layer systems generally require 10-25% Cr, and often use 5-15% B in order to isolate the magnetic grains in the magnetic layer and reduce noise.\nThere exists a continuing need for high areal density longitudinal magnetic recording media exhibiting high Hcr and high SMNR while overcoming the deficiencies of the prior art solutions."}
{"text": "Speech is a natural form of human expression. Irrespective of age, whether it's a small child or a grown up individual, speech comes naturally to humans and is “the” mode of human interaction. Realizing this preferred mode of human interaction, computer scientists have extended this natural form of human interaction to computing systems, and over the years speech recognition systems have evolved considerably in recognizing human speech. Therefore, whether it's a simple command application or a dictation system, speech recognition systems have become quite useful in human-computer interaction. However, one problem with any speech system is spurious speech recognition, especially in a noisy environment, such as, traffic or a factory establishment. Spurious or background noise severely impacts the ideal use of a speech recognition system and creates difficulties during interaction. This affects the system's reliability since the speech recognition system is unable to distinguish between a speech input and a spurious input."}
{"text": "1. Statement of the Technical Field\nThe present invention relates to the field of markup language processing, and more particularly to the processing of frames in markup language.\n2. Description of the Related Art\nConventional markup can be visually presented through use of a content browser. Content browsers process display attributes embedded in markup to properly format content also contained within the markup. Notable variants of the content browser include the venerable Web browser, as well as the more recent extensible markup language (XML) browser. Regardless of the type of browser, all conventional markup processors are preconfigured to parse and interpret attribute tags embedded in markup. Examples of attribute tags include the well-known hypertext markup language (HTML) tags, <HEAD>, <BODY>, <H1>, <P>, <HREF>, <HTML> and <FRAMESET>.\nIn regard specifically to the <FRAMESET> attribute tag defining a set of displayable frames, many content distributors have incorporated frames in the design of Web pages, as frames allow Web sites to organize the presentation of disparate information in a logical and unified manner. Often, the enablement of multiple, scrollable regions of a single Web page forms the basis for the use of frames within the Web page. For example, it is known to use one frame to present a menu, while a second, adjacent frame can display the content associated with a particular menu choice. As another example, the results of a query can be presented in a first, scrollable frame, while the content associated with a particular result can be presented in an adjacent frame.\nWhile frames can provide the benefit of screen organization and scrolling, frames can pose some difficulty in the support of bidirectional environments. Bidirectional environments relate to the directional manner in which information can be presented in an electronic document. Specifically, western languages such as English present information visually from left to right. In contrast, middle-eastern languages such as Hebrew and Arabic present information visually from right to left. In the course of undertaking internationalizing the presentation interface of computer software, solutions have been proposed which address the problem of bidirectional text. For instance, in one solution, text can be rearranged in a display window according to the underlying language associated with the character codes of each word in the text.\nOther internationalization solutions address not only the underlying directional orientation of text, but also the arrangement of user interface controls in a graphical user interface. For example, in Kaplan, Internationalization with Visual Basic, (SAMS 199x), check boxes, labels, option buttons and text boxes can be horizontally “flipped” in terms of orientation to accommodate a bidirectional environment. The particular user interface control elements which can be flipped in the Kaplan reference, however, are flipped inasmuch as those control elements usually accompany text—hence the need to re-orient the control. By comparison, the teachings of Kaplan explicitly inhibit the re-orientation of other user interface controls such as combo boxes, command buttons, scroll bars, list boxes, picture boxes and, most importantly, frames.\nWhile the use of frames in a conventional GUI as described in the Kaplan reference rightfully inhibits the re-orientation of frame elements because frames play little role in the GUI of a stand-alone application, the same cannot be said of an application whose interface relies upon the presentation attributes of a markup language. Specifically, frames play a crucial role in the presentation interface of a markup language defined user interface. Thus, in the context of markup, the directional implication of a language is not merely limited to the characters which form a word, or the words which form a sentence. Rather, the directional implication of a language can include the layout of the frames within the document itself. More particularly, in a right to left orientation, it can be preferable to horizontally re-orient adjacent frames in a content browser to accommodate a right-to-left environment.\nUnfortunately, existing content browsers and markup languages do not account for bidirectional documents. In fact, in the HTML specification, while the “dir” attribute can specify a directional orientation, including “rtl” and “ltr”—right to left and left to right, respectively—the HTML specification cannot account for a document whose orientation can vary from right to left and left to right. Rather, to accommodate the bidirectional circumstance, separate markup must be maintained for both cases of left to right and right to left configurations. As one skilled in the art will recognize, however, maintaining two sets of markup to support the presentation the same content in different directional orientations requires the maintenance and synchronization of both sets of markup—a distinctly undesirable solution."}
{"text": "1. Field of the invention\nThe present invention relates to landscape edging used to demarcate areas of lawns, gardens and the like, and more particularly to an illuminated landscape edging.\n2. Description of the Prior Art\nLandscape edging is very widely used to provide demarcation between areas of lawns, gardens and the like. For instance, landscape edging may be used to form a barrier between botanical types, wherein, for example, on one side thereof is a lawn and on the other side thereof is a flower garden. In order to provide a barrier to passage therethrough of flora, landscape edging is implanted into the ground along the line of selected demarcation, and a portion thereof is located above the surface of the ground.\nOne very common type of landscape edging is depicted in FIG. 1. This prior art landscape edging 10 is characterized by a planar member 12 which is structured to be implanted into the ground and by a tubular member 14 which is structured to be located above ground. The planar member 12 is provided with a medial ribbing 16 and usually provided with an end hook 18. Both the medial ribbing 16 and the end hook 18 serve to anchor the planar member into the ground. This prior art landscape edging 10 is constructed of a plastic material, which is believed to be invariably black and opaque. The prior art landscape edging is elongated and is supplied in long rolls which may be cut into sections to suit the length needs of a particular landscaping job. Sections of the prior art landscape edging are serially joined by a tubular connector 20 inserting into open ends of adjoining tubular members 14.\nIncreasingly, homeowners are installing landscape lighting to illuminate their lawns and other landscape features. Usually, this form of illumination is in the form of a series of discrete lighting units which are partly implanted into the ground and which operate on a low voltage line supplied from the home utility service. These units can sometimes be the target of vandals or may be the accidental victims of an untoward incident, such as being hit by a lawn mower.\nWhat is needed in the art is to somehow combine the benefits of landscape edging with the attractiveness of landscape lighting."}
{"text": "1. Field of the Invention\nThe present invention relates to a laminated gas sensor element, a gas sensor equipped with the laminated gas sensor element, and a method for manufacturing the laminated gas sensor element.\n2. Description of the Related Art\nConventionally, known gas sensors include those for detecting a particular gas component contained in exhaust gas or for measuring the concentration of the particular gas component. Some of these gas sensors employ a laminated gas sensor element configured such that a pair of electrodes are provided on a solid electrolyte body. The solid electrolyte body of zirconia or the like used in the gas sensor element becomes active at a high temperature of 300° C. or higher. Usually, the gas sensor element is used in an activated state upon heating by a heater laminated on the solid electrolyte body. In this case, adhesion, to the gas sensor element, of a water or oil droplet contained in a gas to be measured (hereinafter, also referred to as adhesion of water) may generate a crack in the gas sensor element due to thermal shock. According to a known technique for solving this problem, a distal end portion (detecting portion) of the gas sensor element which is exposed to a gas to be measured is protected with a porous protection layer. For example, Patent Document 1 discloses a technique for restraining the generation of a crack at corner portions of a gas sensor element by increasing the thickness of a protection layer at the corner portions, which are susceptible to the generation of a crack.\n[Patent Document 1] Japanese Patent Application Laid-Open (kokai) No. 2003-322632\n3. Problems to be Solved by the Invention\nHowever, even in the case of a laminated gas sensor element whose detecting portion is protected with a porous protection layer as in the case of Patent Document 1, the protection layer may fail to provide sufficient protection. Namely, adhesion thereto of a larger amount of water droplets or oil droplets contained in a gas to be measured (adhesion of water), potentially results in the generation of a crack in the laminated gas sensor element due to thermal shock."}
{"text": "The present invention relates to a new and distinct cultivar of Begonia plant, botanically known as Begonia hybrida, commercially referred to as a Rex Begonia and hereinafter referred to by the name ‘KRBELSS01’.\nThe new Begonia plant is a product of a planned breeding program conducted by the Inventor in Ermelo, The Netherlands. The objective of the breeding program is to create new Rex Begonia plants with unique and attractive leaf coloration.\nThe new Begonia plant originated from a cross-pollination made by the Inventor in Ermelo, The Netherlands in January, 2009 of a proprietary selection of Begonia hybrida identified as code number KV07B0399-003, not patented, as the female, or seed, parent with a proprietary selection of Begonia hybrida identified as code number KV06B0370-002, not patented, as the male, or pollen, parent. The new Begonia plant was discovered and selected by the Inventor as a single flowering plant from within the progeny of the stated cross-pollination in a controlled greenhouse environment in Ermelo, The Netherlands in January, 2010. Asexual reproduction of the new Begonia plant by leaf cuttings in a controlled greenhouse environment in Ermelo, The Netherlands since February, 2010 has shown that the unique features of this new Begonia plant are stable and reproduced true to type in successive generations."}
{"text": "Fluid flow pipe systems provide the infrastructure necessary to accommodate the flow of fluids, such as, liquids, gases, or slurries, etc, in residential and commercial applications. For example, such systems are commonly found in residential homes, office buildings, factories, car washes, swimming pools, irrigation installations, septic installations, etc.\nRegardless of the applications, most fluid flow pipe systems share certain common elements, such as:                (a) pipes: pipes are hollow cylindrical members configured to conduct or transfer fluids. Pipes may comprise metal materials (e.g., carbon steel, stainless steel, copper, copper alloys, copper-nickel, etc.) or plastic materials (e.g., PVC, CPVC, UPVC, FRP, FRE PEX, etc.);        (b) pipe fittings: pipe fittings are attachments that are placed or affixed to the end pipe ends to provide versatility in changing the direction of fluid flow, distribution, increase or decrease flow capacity and interconnections. Pipe fittings may comprise a variety of configurations (e.g., elbows, bends, returns, tees, crosses, reducers, end caps, plugs, nipples, unions, couplings, bosses, etc.);        (c) valves: valves are mechanical devices configured to control the flow and pressure within the pipe system (e.g., gates, globes, plugs, balls, butterflies, checks, diaphragms, pinches, pressure reliefs, control valves, etc.;        (d) fixtures: fixtures are elements or amenities disposed at the end of pipe runs (e.g., faucets, toilets, showers, bidets, spigots, sprinklers, icemakers, etc.); and        (e) pipe system service equipment: such equipment includes, but not limited to, water heaters, filtration units, desalinization units, water softeners, back-flow preventers, sump-pump units, etc.        \nFor a variety of reasons, including but not limited to, environmental conditions, rust/corrosion, pressure variations, component fatigue/failure, substandard installation practices, etc., it is not uncommon for elements of fluid flow pipe systems to deteriorate or become defective, thereby requiring repair or replacement servicing. Repair/replacement servicing often requires that the fluid in the locally-affected area or defective equipment, or sometimes in the entire pipe system infrastructure, be sufficiently evacuated prior to the repair or replacement, resulting in time-consuming delays, inefficiencies, and customer inconvenience."}
{"text": "1. Field of the Invention\nThe invention relates to degaussing coils or degaussers designed to be fitted into cathode-ray tubes for color televisions. The degaussing coil is constituted by a winding of enamelled electrical wire shielded by a heat-sealed insulator sheath. This wire is fastened to the glass fixtures of the mask picture tube and it is connected to the electrical circuit of the television set by a pair of conductors that end in a bipolar connector. The conductors are electrically connected to the two ends of the winding in a region that is not covered by the insulator sheath. This uncovered region is then confined in an insulating envelope that ensures the electrical shielding and the mechanical strength of the connected parts. The envelope can be formed by taping or by two shells of synthetic organic resin forming an interconnection case.\n2. Discussion of the Background\nCoils manufactured according to taping techniques or by the positioning of a case require equipment whose difficulty of handling entail considerable extra costs. These coils may have a defect in the sealing quality at the interconnection and they lack aesthetic quality."}
{"text": "A lens array optical system including a plurality of lens portions is used for an imaging lens for a compound-eye imaging apparatus, a secondary image forming lens in an autofocus module of a single-lens reflex camera, an illumination lens in a semiconductor exposure apparatus, a condenser lens in a liquid crystal projector panel, and the like. The lens array optical system can be manufactured at low cost with optical plastic having high processability. However, the optical plastic not only has a high thermal expansion coefficient bus also is inferior in durability and light transparency, and accordingly is not suitable to be used stably on a high-temperature or high-humidity severe condition. Therefore, optical glass that has a low thermal expansion coefficient and is superior in durability under high temperature or high humidity is required to be used for the lens array optical system in order to enable stable use in various environments.\nAs a method for manufacturing a lens array optical system made of optical glass having multiple spherical or aspheric lenses, a method is proposed in which preforms made of glass with a larger curvature than a concave curvature of a forming die are placed, one by one, in transfer portions and are hot stamped (refer to Patent Literature 1). An excess of the preforms over the volumes of cavities for forming optical surfaces of the multiple lenses flows into junctions of the cavities and are integrally fused. However, there is a problem with the method of Patent Literature 1 that it takes much time and trouble to place the preforms on the optical transfer surfaces on the die. Moreover, air may remain in the junction (fused portion or joint) due to variations in the preforms to result in a reduction in the strength of the junction of the preforms. Moreover, there is also a problem that as the number of lens portions increases, the forming process becomes more difficult, and accordingly it becomes impossible to obtain a lens array having an intended optical surface shape.\nMoreover, as another method for manufacturing a lens array optical system, a method is also proposed which does not use a die, supplies a liquid material to a substrate where a plurality of through-holes is formed, and forms a plurality of lenses on the substrate (refer to Patent Literature 2). However, there is a problem with the method of Patent Literature 2 that it is difficult to form an optical surface of the lens into an aspheric shape so that it is not possible to appropriately design the optical surface of the lens portion and optically correct various aberrations.\nMoreover, as another method for manufacturing an optical glass lens array optical system, the drop method is known which forms a lens by dropping molten glass into a die. The drop method has few constraints on the shape of an optical surface since the forming process is performed in a state where the viscosity of the glass is relatively low. Moreover, a reduction in the occurrence of forming failure ascribable to the fluidity of a material can be expected. As the method for manufacturing a lens array by the drop method, Patent Literature 3 describes a method for manufacturing a lens array having 2×2 lens portions by dropping a glass drop on the center of one of dies having 2×2 optical transfer surfaces and pressing it with the other die. However, Patent Literature 3 does not describe the production of a lens array having lens portions more than 2×2, such as 4×4 or 5×5, in a lattice form. If an attempt is made to obtain a lens array having 4×4 lens portions by dropping a glass drop on a center portion of, for example, 4×4 optical transfer surfaces in accordance with the manufacturing method described in Patent Literature 3, an optical transfer surface that has been insufficiently filled with the glass may be produced as described below. Moreover, in Patent Literature 3, a projection is provided to change the flow of the glass drop to make adjustments such that the glass drop flows along the transfer surface of an optical surface near its edge side close to the projection, the edge side having a large inclination angle. Accordingly, an attempt is made to transfer the optical surface with high accuracy. However, there are problems that the manufacturing cost may increase if a glass projection is provided to the die as a preliminary process, and the die processing cost is high and the life of the die is influenced if the projection is formed upon die processing. If an attempt is made to drop a plurality of glass drops individually and respectively onto the optical transfer surfaces, forming failure may be invited, or gas entrainment may occur upon reaching the die, due to time differences in die arrival timings among the glass drops. Moreover, a problem also arises in which it becomes difficult to perform a forming process by pressing to a desired thickness since it becomes easy to become cold and harden as the result of a reduction in the volumes of the individual glass drops."}
{"text": "A known voltage converter and a known inverter have been used when driving a three-phase motor. The voltage converter boosts an input voltage to a predetermined voltage. The inverter changes a frequency of an output of the voltage converter. Meanwhile, plural switching elements are provided at the voltage converter and the inverter. The switching elements self-heat by energization, and electrical characteristics and a rating are specified in response to an ambient temperature level of the switching elements. Here, the ambient temperature is desired to be considered upon the use of the switching elements. The technology in which the switching elements are used under the consideration of the ambient temperature is disclosed in JP2013-48515A (hereinafter referred to as Patent reference 1).\nAn electric automobile disclosed in Patent reference 1 is provided with an inverter supplying electric power to a three-phase motor. The inverter includes plural switching elements. The electric automobile is provided with a temperature sensor, a current sensor, and a voltage sensor. The temperature sensor measures a temperature level of a refrigerant cooling the switching elements. The current sensor measures an output current of the inverter. The voltage sensor measures an input voltage inputted to the inverter. A temperature correction value is calculated based on a measurement data of the current sensor per switching element, a measurement data of the voltage sensor per switching element, and a duty ratio of the switching element. The temperature correction value is added to a measurement result of the temperature level of the refrigerant to estimate the temperature level of the switching elements.\nAccording to a device disclosed in Patent reference 1, detection results of the plural sensors are used to estimate the temperature level of the switching elements. In this case, because the detection results include errors of measurement, the errors of measurement included in the estimate results of the temperature levels of the switching elements may be increased. Accordingly, when the switching elements are energized, the electric characteristics and the rating are desired to include margins. Accordingly, the device disclosed in Patent reference 1 may not be able to fully use a capability of the switching elements.\nA need thus exists for an energization control system and a sensor unit which is not susceptible to the drawback mentioned above."}
{"text": "Competitive electronic sports, or esports, allow players to participate in virtual or digital environments. Traditional sports competition occurs in a physical game space, such as in a football stadium or in a basketball arena. Thus, in traditional sports competition, multiple points of view of the game play may be provided to viewers by placing cameras at different positions within the physical space. However, with esports, as the competition occurs in a virtual or digital environment, providing viewers with multiple views of the competition action is not as simple as placing multiple cameras within a physical space. Therefore, providing a viewer with multiple points of view of the competition action comprises providing the viewer with multiple video streams, each from a different point of view. This method of providing a viewer with multiple points of view of the competition action may require a significant amount of bandwidth and processing resources. For example, if the multiple points of view include more than, for example, four points of view, the viewer may not have access to sufficient bandwidth to view the multiple points of view. And as will be recognized, an esports competition may comprise two teams of five players competing within a virtual or digital environment and a viewer may wish to watch the competition action from the view point of each of the players. This would require the viewer to have access to sufficient bandwidth for viewing ten different video streams.\nTherefore, a need exists for improved methods, apparatus, systems, computer program products, computing devices, computing entities, and/or the like for providing a viewer with multiple viewpoints of an event occurring in a virtual or digital environment that does not require an excessive amount of bandwidth."}
{"text": "There are numerous applications in which it is desirable to be able to detect underwater objects from stand-off distances or remote locations. Some conventional approaches to the detection of underwater objects have included the use of unmanned underwater vehicles (UUVs) that search for objects within the water, and aerial vehicles (such as helicopters) that can drag nets to comb the waters for near-surface objects. Other approaches have involved the use of remote sensors, optionally deployed on buoys, or other relatively-static deployed structures on the surface of the water, that look for objects or indirect indicators of the presence of an object of interest (e.g., a vessel wake); however these methods tend to suffer from low confidence-of-detection, high false alarm rates, and/or low area coverage rates."}
{"text": "For the polymerization of olefins the Ziegler-Natta-catalyst system is generally used, which consists of a so-called procatalyst and a cocatalyst. The procatalyst is based on a compound of a transition metal belonging to any of the groups IV(A) to VIII(A) (Hubbard) of the periodical system of elements and the cocatalyst is based on an organometallic compound of a metal belonging to any of the groups I (A) to III (A) (Hubbard) of the periodical system of elements.\nNowadays, the procatalysts typically comprise an inert carrier, on which the actual active catalyst component i.e. the transition metal compound or the mixture of the complex formed by the catalytical compounds has been layered. The morphology and size distribution of the particles of such a carrier are most significant for the activity of the catalyst and the properties of the polymer obtained through the catalyst. With an active catalyst you are, namely, able to produce a polymer, from which, due to its purity, no catalyst residues need to be removed. The morphology of the carrier, again, influences the morphology of the polymer product itself for it has been noticed that the morphology of the catalyst is repeated in the structure of the polymer (the so-called replica phenomenon). When a flowing product polymer having the desired morphology and a narrow particular size distribution is wanted, which is desirable in view of the many aims of use of the processing processes, the properties of the carrier shall by the aid of the replica phenomenon be made similar.\nNowadays, the Ziegler-Natta type procatalysts typically comprise a magnesium based carrier, such as magnesium chloride, which has been treated with a transition metal compound of titanium halide, such as titanium tetrachloride, and sometimes also with an electron donor compound. It is also known that the carrier can be brought to an advantageous and equal-size crystal form by letting it crystallize as a complex from any of its crystal solvents.\nThe treatment of such a carrier with a crystal solvent has been disclosed among others in the U.S. Pat. No. 4,071,674, in which the procatalyst based on a transition metal compound has been prepared by bringing a titanium or vanadium compound to react with a reaction product that has been formed when magnesium dihalide and the addition product of alcohol are reacting with an organometallic compound of a metal of any of the groups I-III. The preparation of the procatalyst begins with the addition of alcohol dropwise to the suspension of magnesium dihalide, after which the organometallic compound is added dropwise to the reaction mixture. After agitation the preactivated carrier is activated by adding titanium tetrachloride to the mixture. The adding stages of this kind of a method are primitive and do not at all allow regulation of the morphology of the procatalyst in the way desired.\nTreatment with a crystal solvent has also been described in patent application JP-59-215301. In this publication the carrier complex (10 g of MgCl.sub.2 and 24.2 g of EtOH) have been prepared by an emulsifying technique. The carrier complex melt has been dispersed into n-decane as spheroidal melt particles. Thereafter the carrier particles in the emulsion have been shock coagulated by transferring the emulsion into cold hydrocarbon medium. A drawback of this method is, among others, that such components are needed in the preparation of the carrier that are not useful at later stages of the catalyst preparation and this presupposes the existance of a purifying and recirculation equipment for this purpose.\nThe patent family comprising, among others, the EP publication 65700 and the U.S. Pat. No. 4,421,674 which claims priority from the Italian application IT 2,188,181 (810521), deal with a method for the preparation of a catalyst, which is particularly active in the polymerization of gaseous ethylene.\nIn the process titanium halide is brought to react with a magnesium chloride catalyst carrier being in the form of microballs, after which the reaction product particles are recovered by physical means and are mixed together with an organometallic cocatalyst product.\nCharacterizing of this method representing the prior art is that:\nAccording to FI patent publication 80055 (Neste Oy) the above-mentioned carrier complex formed by the carrier and the crystal solvent can be melted to clear liquid. When this kind of a liquid is conducted through the spray nozzle to the spray space cooled with cold nitrogen gas it crystallizes to spheroidal small particles of the carrier complex which are very flowable and loose. The process took in practice place so that MgCl.sub.2 and C.sub.2 H.sub.5 OH were melted at the temperature 110.degree. to 130.degree. C. to a clear melt. Then the clear homogenized mixture was fed through a nozzle dispersing it into a cooled spray chamber. The spraying gas used in the spraying was dry nitrogen having the temperature +130.degree. C. and as cooling medium dry nitrogen was conducted to the spraying chamber, the temperature of which was 20.degree. C. As the nozzle a gas fluidizing nozzle was used.\nThis kind of a method produces very flowable and loose particles. Moreover, the carrier complex crystallizes without any evaporation of the crystal solvent. When such a carrier is brought into contact with titanium compound, abundantly for catalytically active complexes between MgCl.sub.2 and the titanium compound are formed on the surface of the carrier when the crystal solvent leaves.\nAccordingly, two methods of preparation based on the spraying are prior known in the field. The spray drying method patents are based on a fairly complete drying of the carrier liquid from ethanol (C.sub.2 H.sub.5 OH) after the spraying. Hereby, the carrier has usually been dried at a temperature exceeding 150.degree. C., whereby a major portion of the alcohol of the complex is evaporated. In the spray-drying method a carrier product is typically obtained, the alcohol concentration of which is between 15 to 25% by weight and anyway below 30% by weight.\nOne central drawback of spray-drying is the bad morphology of the carrier obtained and the wide particle size distribution, which is caused by the breaking of the particles during the process. In FIGS. 1 and 2 an electron microscope picture of the carrier is presented. The close-up picture of FIG. 2 shows that the spheroidal particles are either compressed or then they have been broken into crust fragments of the hollow ball. On the basis of the physical process and the particles created in it, it can be supposed that the removal of C.sub.2 H.sub.5 OH at temperatures exceeding +150.degree. C. results in the formation of a gas state inside the particles, which then results in the disintegration of the spheroidal particles into crust fragments or their becoming compressed when the gas cools down in the hollow space. The spray-drying will, anyhow, result in a very poor carrier morphology.\nIt can thus be said that while the spraying steps of prior spray-drying methods were easy to carry out and their yield of produced carrier is good, the activity of the procatalyst obtained by activation with TiCl.sub.4 when polymerizing PP (polypropylene) is only between 10 to 12 kg of PP/g of catalyst and the morphology of the polymer obtained by it is very poor producing even up to 60 to 70% by weight of finely-divided material (d&lt;1 mm), the bulk density being of the order of about 0.2 g/ml. The above-mentioned activity of the procatalyst obtained by spray-drying is lower than that of the procatalysts obtained by other methods, which is due to the fact that in spray-drying the major portion of the C.sub.2 H.sub.5 OH participating in the TiCl.sub.4 activation leaves the carrier before it is reacted with TiCl.sub.4.\nThe spray-crystallization of such a melt MgCl.sub.2 --C.sub.2 H.sub.5 OH complex that contains on average 3.5 or more C.sub.2 H.sub.5 OH molecules as per each MgCl.sub.2 molecule is easy and gives a relatively good yield, whereby the result is a procatalyst having a good activity i.e. of the order of about 15 kg of PP/g of catalyst. It has, however, been noticed that the C.sub.2 H.sub.5 OH amount of the carrier is the same as the respective amount of the melt and too high causing during the TiCl.sub.4 activation heavy formation of HCl gas, which breaks the procatalyst particles and results in a poor procatalyst morphology and a poor particle size distribution. In the polymerization polypropylene having a poor morphology and about 25 to 55% by weight of finely-divided substance (d&lt;1 mm) is obtained. The bulk density of the polymer obtained is 0.40 to 0.44 g/ml.\nIn the spray-crystallization of melt MgCl.sub.2 --C.sub.2 H.sub.5 OH complex which contains less i.e. in the average about 2.9 molecules of C.sub.2 H.sub.5 OH as per each MgCl.sub.2 molecule, big difficulties are again encountered when spraying melt into the crystallization chamber. Melt with little alcohol apparently is too viscous and causes clogging of the nozzle and formation of too big particles e.g. as a result of solidification. The carrier yield is very low remaining below 10% in the polymerization of propylene a good activity of the procatalyst is obtained i.e. about 15 kg of PP/g of catalyst, but the morphology of the polypropylene is poor and it contains 20 to 40% of finely-divided material (d&lt;1 mm), the bulk density being between 0.40 to 0.44.\nAccordingly, it seems that the spray-crystallization results in a procatalyst whose activity is medium, and whose polymer products contain a very great portion of finely-divided material and possess a polymer morphology which is poor."}
{"text": "1. Technical Field\nThe present disclosure relates generally to integrated circuit (IC) design. More particularly, and not by way of any limitation, the present disclosure is directed to an architecture and associated system and method for compressing repair data in an IC design having a plurality of memory instances.\n2. Description of Related Art\nModern integrated circuits (ICs) such as application-specific ICs or ASICs may contain a large number of repairable memories. It is not uncommon to have 1000-2000 memories in some ASIC implementations. In such cases storing reconfiguration data to repair the memories may require upwards of 10,000 to 20,000 bits. Usually, this data may be split into multiple containers (e.g., banks of fuses) across the chip area. However, in some cases allocation of such a large container on an IC may not be possible."}
{"text": "This application claims benefit of priority under 35 U.S.C. xc2xa7120 to International Application PCT/EP99/03219, filed May 11, 1999, and under 35 U.S.C. xc2xa7119 to DE19821009.4, filed May 11, 1993.\nThe present invention relates to a process for extracting carotenes from carotene-containing materials, in particular from fats and oils of biological origin. though the invention shall chiefly be explained through the example of carotene, the term xe2x80x9ccarotenesxe2x80x9d encompasses not only the isomers of carotene but also carotenoids.\nAccording to a definition by IUPAC, carotenoids are chemical compounds of aliphatic or aliphatic-alicyclic structure having conjugated double bonds and comprising 3 to 8 (or even more) isoprene moieties.\nThe carotenoids are the most important group of natural colouring materials most widely occurring in plants and animals. They are fat-soluble, nitrogen-free, yellow to violet materials, wherein the uninterrupted sequence of isoprene moieties and thus the chromophoric accumulation of conjugated double bonds gives rise to colouring. All carotenoids being polyenes, they exhibit a blue solution color in concentrated sulphuric acid.\nCarotene is the carotenoid which has been known for the longest time. It was isolated for the first time by H. Wackenroder (in 1831) from carrots. About 100 years after its discovery, Kuhn found in 1931 that the naturally occurring carotene is composed of three isomers which he named xcex1 carotene, xcex2-carotene and xcex3-carotene. In the untreated carotene of carrots, they are contained at proportions of about 15, 85, and 0.1%, respectively. Ever since the first carotenoid syntheses by P. Karrer and H. H. Inhoffen (1950), a rapid development commenced in the field of commercial syntheses of carotenoids.\nAll three carotene isomers possess the same fundamental structure, comprising one xcex2-ionone ring structure at one molecule end, 9 conjugated double bonds and 8 branchings. They differ only in the structure of the other molecule end.\nDepending on its origin, carotene is a varying mixture of the structurally isomeric polyene hydrocarbons C40H56: all-trans-xcex1-carotene, all-trans-xcex2-carotene, all-trans-xcex3-carotene. Depending on processing and starting material, cis isomers may also occur.\nCarotene occurs not only in carrots but also in numerous other plants, and particularly accompanying chlorophyll it counts among the most frequently occurring plant coloring materials. Carotenes are present in vegetable oils at moreor less elevated concentrations. The carotene contents are particularly high in palm oil, in the ranging from 0.05 to 0.2%.\nRaw carotene is a dark coppery to cinnabar-colored crystal powder of wax-like consistency. Its solution color is yellow to orange.\nxcex2-Carotene is the carotene occurring most frequently in animals and vegetables. It forms deeply violet prisms (crystallised from benzene/methanol) or polyhedra (from petroleum ether). Its solubility in grams per 100 cm3 of solvent at 19xc2x0 C. is 5.5 in carbon disulfide; 0.35 in benzene; 0.1 in petroleum ether; 0.008 in ethanol. It is optically inactive. xcex1-Carotene forms deeply violet prisms (from benzene/methanol) or polyhedra (from petroleum ether). It is optically active and more readily soluble than xcex2-carotene. xcex3-Carotene forms dark red prisms (from benzene/methanol). It is optically inactive and also more readily soluble than xcex2-carotene.\nCarotene is merchandised in the form of the following preparations: crystallised xcex2-carotene in vacuum ampoules; technically crystallised at about 80%; carotene concentrates in powder form at about 3%; carotene solutions in vegetable oils at about 0.3 to 0.5%.\nCarotenes, being polyene hydrocarbons, exhibit good solubility in tetrahydrofuran, carbon disulfide, benzene, chloroform and oils, low solubility in ether and petroleum ether. All carotenes are insoluble in water. In fats and oils of biological origin, their solubility typically is 2 to 3%. Almost all carotenes have high melting points (e.g., xcex2-carotene 183xc2x0 C.). They are sensitive to acids, oxygen and exposure to light.\nThe carotenes, have gained importance in nutrient technology as antioxidants and as natural colouring materials. They predominantly serve for coloring fats and oils, for vitamin enrichment of margarine, nutrient preparations, and pharmaceuticals, as an addition to concentrated feed in rearing young animals and to ice creams or sherbets and milk preparations.\nThe carotenoids, being hydrocarbons or compounds closely related to hydrocarbons, are markedly lipophilic. (Carotenoids of the polyene alcohol type dissolve well in alcohol and acetone.) All-trans forms do not dissolve as well as corresponding cis compounds. Carotenoids occur in nature almost exclusively in solution in the lipoid particles of the cell, for which reason they are at times referred to as lipochrome colouring materials. From the dried cell substance they can be extracted only by means of lipophilic solvents. The commonly known fact that the carotenoids of our nutrient plants, e.g., the lycopene of the tomato, enrich in the fat globules of the foods prepared from them, is founded in this lipophilic character.\nH. V. Euler realised (in 1928) the property of carotene being provitamin A. xcex2-Carotene is related to vitamin A, which latter one may be conceived to be formed by cleavage of the central Cxe2x95x90C double bond of the carotene with concomitant incorporation of two molecules of water. Biosynthesis of vitamin A also appears to unfold in this way. At least the animal body is capable of transforming in most cases carotene, especially xcex2-carotene, into vitamin A. This is of interest inasmuch as a vitamin A deficiency may be remedied through easily obtained foods which contain carotene.\nFor obtaining the carotene, the natural material is dried mostly at temperatures below 50xc2x0 C., and the carotenes are extracted with lipophilic solvents. Jointly dissolved accompanying substances are separated out either by saponification or by freezing out. Crystallised raw carotene (75-90%) or oil-based carotene concentrates (1-20%) are obtained. Of the isomers, only xcex2-carotene is obtained pure. Starting products are either such having a high carotene content, like carrots (1 g of raw carotene/kg of dry substance), raw palm oil (up to 3 g/kg), pumpkin seeds or green plants such as lucerne (alfalfa), stinging nettle or broccoli (Brassica), from which xanthophyll and chlorophyll are furthermore obtained besides carotenes.\nFor the extraction of carotenes from dry plants, petroleum ether, naphtha, less frequently benzene, fatty oils are used, or also trichloroethylene in the case of red palm oil. For obtaining them from carrots, autumn carrots [Herbstmxc3x6hren] are suited best, which should contain at least 10 mg % [mg/100 ml] of carotene for 10-15% dry substance. One starts out from dried carrots (dry processing method) or from material processed in the autoclave and pressed hydraulically (wet processing method). The dry method procures a yield of 67 g of carotene/1000 kg of carrots. In the wet method, extraction is performed with tetrahydrofuran and provides a yield of 62-72 g of carotene from 1000 kg of carrots having a carotene content of 90-100 g.\nThe most important cultivation areas for oil palms are Africa, Indonesia, Malaysia and Brazil. Depending on their origin and freshness condition, palm oils have a bright yellow (predominantly xcex1-carotene), red (lycopene), orange (predominantly xcex2-carotene) or reddish brown (presence of chlorophyll) coloration. In general, the raw oil is processed with alkali into low-water soaps from which the carotenoids are extracted with benzene, trichloroethylene or petroleum ether. The yield is approximately 80-90% of the overall carotenoids of the oil. From 1000 kg of palm oil, one obtains about 3-4 kg of carotene extract (containing 20% carotene, 6.5% lycopene and 6% volatile oils).\nThe recovery of the carotenes present in oils and fats by way of saponification of the glycerides and subsequent extraction of the carotenes from the soaps with lipophilic solvents is a costly process. Moreover carotenes from fats and oils, which pass into the food sector without being chemically modified, are not accessible in this way.\nDE-A-195 31 245 discloses a process for the extraction of carotenes, in particular xcex2-carotene, from solid biological materials by means of a liquid organic extractant mixture consisting of ethyl acetate and/or butyl acetate and 3 to wt. % of an oil of biological origin, relative to the ethyl acetate and/or butyl acetate, at temperatures between 40 and 125xc2x0 C.\nDE-A-44 29 506 discloses a process for the extraction of carotenoids from natural starting materials by means of dense gases, wherein the pre-dried natural starting material is extracted with dense propane and/or dense butane, optionally in the presence of an organic entrainer, at temperature between 20 and 100xc2x0 C. and pressures between 10 to 200 bar.\nEP-A-0 455 425 discloses a process for the preparation of concentrates of coloring agents, wherein said concentrates of natural coloring agents such as carotene are prepared from organic media, particularly from palm oil, by a process in which the oil, together with a volatile solvent, is subjected to gel permeation chromatography.\nEP-A-0 242 148 discloses a method for purification of a carotene-containing concentrate by chromatography in which a liquid containing a carotene-containing concentrate is fed to a filler to adsorb the concentrate on the filler, and the carotene component is subsequently eluted and collected by feeding an eluant to the filler.\nEP-A-0 239 949 discloses a process for the manufacture of powdery, finely divided preparations of carotenoids in which the carotenoid essentially has a particle size of less than 0.5 xcexcm, comprising dissolving the carotenoid in a volatile, water-miscible organic solvent at temperatures between 50 and 240xc2x0 C., under atmospheric or superatmospheric pressure, in less than 10 seconds, and precipitating the dissolved carotenoid in colloidal disperse form by rapid mixing with milk at from 0 to 50xc2x0 C.\nWO-A-98/03480 discloses a process for recovering highly pure xcex2-carotene crystals from a crude crystal preparation obtained from a natural source. To remove impurities, the crude crystals are stirred in a solvent in which xcex2-carotene has a low solubility, whereupon the crystals are filtered off and washed with fresh solvent. Using this process, a natural crystalline xcex2-carotene preparation with a very high purity is obtained.\nWO-A-96/29306 discloses a process for the recovery of carotene from native fat or oil, in particular from palm oil, wherein the native fat or oil is converted with an alkanol with a C chain length of up to 4 catalytically into fatty alkyl ester and glycerine in a prior art manner. The ester phase of the reaction mixture is subjected to distillation to separate the fatty acid alkyl ester. The distillation residue obtained in the second step is saponified, carotene is extracted from the product obtained in the third step and the extract phase is concentrated by evaporation. A yield of at least some 80% is obtained. At the same time a fatty acid alkyl ester is provided for further processing into fatty alcohol.\nJP-A-09048927, Patent Abstract of Japan, discloses a process to provide a high-quality odorless carotenoid pigment used for foods, drugs, quasi-drugs and cosmetics by removing the inherent characteristic odor and contaminants from a material containing a carotenoid pigment, wherein a purified carotenoid pigment is obtained by mixing a material containing 1% or above carotenoid with a water-soluble organic solvent and then with a fat-soluble organic solvent or by mixing the material with a fat-soluble organic solvent and then with a water-soluble organic solvent.\nA method was sought to obtain the carotenes by direct extraction of carotene-containing materials, in particular of fats and oils of biological origin.\nSubject matter of the present invention is a process for extracting carotenes from carotene-containing materials, which comprises extracting a carotene-containing material with an extractant comprising at least one member selected from the group consisting of acetonitrile, N-methylpyrrolidone (NMP), N,N-dimethylformamide (DMF), tetrahydrofuran, N,N-dimethylacetamide, furfurol, morpholine, 4-formylmorpholine, 4-acetylmorpholine, 4-methylmorpholine (NMM), 4-phenylmorpholine, with the formation of two liquid phases, one of which is a carotene-depleted raffinate phase and the other one is a carotene-enriched extract phase; and separating the two liquid phases.\nThere is a large number of carotene-containing materials, in particular fats and oils of biological origin, to which the method according to the invention is applicable and which have economically interesting contents of carotenoids. Palm oil is particularly preferred in this context. Further examples to be named include: soybean oil, rapeseed oil, corn oil, hemp oil, sunflower oil, olive oil, linseed oil, beet oil, castor oil, sesame oil, cocoa butter, rice germ oil, cottonseed oil, grape kernel oil, wheat germ oil, palm kernel oil, coconut oil, peanut oil, safflower oil, and mixtures thereof.\nThe extraction according to the present invention is preferably performed in countercurrent flow. This allows to exploit even small differences in distribution coefficient.\nThe extractant dissolved in the raffinate phase and/or the extract phase may (preferably between room temperature and 50xc2x0 C., however above the melting point of the fat or oil) be removed by liquid-liquid extraction with water. As a result of the very high distribution coefficient, high concentrations of the extractant in the aqueous solution may hereby be obtained. By rectification, the water and the extractant may be separated from each other. The recovered extractant may be recirculated.\nAnother possibility of purifying the raffinate phase and/or extract phase with respect to the extractants consists in removing the extractant from raffinate phase and/or extract phase by stripping with nitrogen, steam (steam distillation) or alcohol vapor, selected from among methanol vapor, ethanol vapor and isopropanol vapor, preferably under reduced pressure.\nMoreover the extractant may be removed from the raffinate phase and/or extract phase by extraction with dense gases, selected from the group consisting of carbon dioxide, propane, propene, butane and mixtures thereof, preferably under near-critical conditions. As an example, 4-formylmorpholine is soluble at 70xc2x0 C. and 120 bar at a concentration of 9 wt. % in carbon dioxide.\nNear-critical conditions, in the framework of this invention, preferably comprise reduced temperatures of from about 0.7 to about 1.3. The term xe2x80x9creduced temperaturexe2x80x9d is defined by the equation: Tr=T/Tc, wherein Tr is the reduced temperature, T is the actual working temperature and Tc is the critical temperature of the respective gas. The unit of T and Tc is Kelvin [K]. Tr does not have a unit, therefore. Butane, for instance, exhibits a reduced temperature of about 0,69 at room temperature.\nThe carotenes enriched in the extract phase may be recovered from said extract phase by crystallization and filtration. Preferably the filtered-off carotenes are commercialised while still oil-wet because pure carotenes have a tendency to undergo violent polymerisation reactions. The carotenes obtained in this way may, however, also be purified in accordance with conventional methods and sealed in vacuum ampoules. The remaining, carotene-saturated oil may be commercialised as such. As an alternative, this oil might also be recycled.\nVarious extractants, such as, for example, dimethylformamide, 4-formylmorpholine, etc., form two-phase systems with most oils and fats without the addition of auxiliary agents.\nWhere required, small quantities of water or of an alkanol entirely miscible with water, preferably methanol, are added to the extractants in order to bring about a two-phase condition.\nAnother option of establishing a two-phase condition in the systems of the extractants and the carotene-containing oil or fat consists in adding carbon dioxide, propane, propene or butane under elevated pressure to the mixture. Herein pressure and temperature are selected such that the binary system of gaseous components and solvent is single-phase, i.e., supercritical.\nA third method for reducing the mutual solubility of extractant and oil is to add an alkane such as, for example, propane, butane, heptane or petroleum ether to the oil. Dimethylformamide and 4-formylmorpholine, for example, are miscible with most organic solvents with the exception of aliphatic hydrocarbons and glycerol triesters.\nThe slimy substances, solids, and free fatty acids are preferably removed in a maximum possible degree by degumming, filtration and deacidification, as well as combinations of these pretreatments. Oxygen should suitably also be removed beforehand.\nThe fat or oil of biological origin, prior to the extraction, may optionally be subjected to a preliminary extraction with a solvent selected from among methanol, ethanol, ethanol/water, isopropanol, isopropanol/water and mixtures thereof. However, this is not a prerequisite for the extraction according to the present invention.\nAs extractants for the extraction according to the invention, mixtures of N-methylpyrrolidone and methanol, preferably containing 30-70 wt. % methanol, in particular 40-60 wt. % methanol, or mixtures of 4-methylmorpholine and water, preferably containing 3-20 wt. % water, in particular 5-15 wt. % water, have been found to be well suited.\nExtractants which are particularly well suited for the extraction according to the present invention are 4-formylmorpholine and mixtures of 4-formylmorpholine and methanol, preferably containing 5-30 wt. % methanol, in particular 5-10 wt. % methanol.\nExtraction in accordance with the invention is preferably performed at temperatures in the range of from about 20 to about 120xc2x0 C., preferably in the range of from about 40 to about 80xc2x0 C. One suitably operates above the melting point of the respective fat or oil."}
{"text": "The present invention relates to the field of electronic design automation and, in particular, to improved techniques for performing optical proximity correction.\nThe manufacture of integrated circuits strives to place ever smaller features onto a given area of the integrated circuit chip. One challenge encountered in this effort to fabricate smaller features is the diffraction of the light used in photolithography. That is, the quality and fidelity of the microlithography stage of very large scale integrated (VLSI) circuit chip production depends on the wavelength of the light source and the size of the features to be printed.\nRecent sub-wavelength lithography approaches aim to use wavelengths that are larger than the minimum feature size to generate the images, (e.g., light with a wavelength of 193 nanometers is being used to generate features with dimensions of 90, 65, or 45 nanometers). This approach, however, requires methods for the correction of degradations and distortions in the final pattern caused by light diffraction. That is, the photolithography mask used to generate the desired circuit pattern includes structures that anticipate and, at least partially correct for, the imperfections arising from striving to fabricate small features.\nA computational simulation of the exposure and lithographic is run and the degradations or distortions are computed with various additions, inclusions and adjustments to the mask design. A mask design is selected that improves the final structure. These methods, commonly known as optical proximity correction (OPC), are mainly dependent on the optical system and mask features and may be computationally intensive. While regions having densely packed features tend to be more prone to distortions (the “proximity” effect), OPC calculations are not limited to such regions and can be advantageously applied to less-dense regions of the circuit.\nOPC typically numerous features in a pattern layout to be computationally processed one or more times. Recent advances in semiconductor manufacturing allow billions of transistors (i.e., multibillion features) to be placed on a single chip. The well-known “Moore's law” postulates that the number of transistors that can be placed on a single chip doubles about every 12-24 months. Unfortunately, despite the advances in the central processing unit (CPU) clock speed and computing power, the gap between the computational power required for OPC calculations and the available CPU processing power keeps increasing. That is, the computing power required to efficiently execute the OPC calculations in a timely manner is growing at a faster rate than the available CPU power in a reasonably priced engineering workstation.\nTo further complicate the issue, the number of masks or layers to which OPC should be applied increases at each new semiconductor device manufacturing node. Since the features are getting smaller with every manufacturing node while the illumination wavelengths remain the same or decrease at a slower rate, the number of neighboring features effecting the fidelity of each feature increases. Therefore, the computational processing power required to perform OPC operations on new chip designs has been increasing at a rate of approximately factors of three or four or more for each successive manufacturing node.\nPresently, the generation of optically corrected masks takes from many hours to several days per mask and the complexity of this process continues to grow. Because the features printed after the OPC process may still be different from the desired features, the impact of each feature on the functionality and performance of the chip needs to be readdressed in an iterative manner. A typical VLSI design process consists of several iterations of mask generation, OPC process, and interpretation of the results. These iterations may contribute several months of delay in the chip qualification and manufacturing process.\nThe persistent time-to-market pressures on new chip designs mandate improved methods to estimate and shorten the impact of the OPC process in the early stages of the design. Since it is computationally prohibitive to perform many iterations of OPC on a full-chip scale, partial or simple model-based OPC approaches are being applied in limited fashion, still necessitating full-chip OPC once the design is completed.\nTherefore, a need exists in the art for improved systems and methods that shorten the time required to perform OPC, improve the accuracy of OPC methods, and are scalable to address larger chip designs."}
{"text": "It is often necessary to perform actuation movements and to generate force effects at different times or consecutively in injection molding machines. The actuation movements occur over more or less large displacements if little or no displacements are to be overcome when force effects are applied. The varying requirements occur in the operation of injection molding machines in connection with, for example, a pusher type screw unit of extruders wherein the plasticized plastics melt is metered into an injection mold by displacing the pusher type screw and the melt is subsequently subjected to a so-called dwell pressure.\nSimilar varying requirements also arise when operating a closing unit of an injection mold where two mold halves are movable towards or away from each other for the respective closing and opening of the injection mold, wherein the injection mold must be subjected to a specific closing force in its closed position.\nIt is therefore an object of the present invention to optimize the operating mode of a true-to-type displacement and/or actuating force driving device for injection molding machines.\nA further object of the invention is to provide an economical displacement and/or actuating force drive device for injection molding machines.\nAnother object of the invention is to provide a device whose electromechanical displacement drive apparatus overcomes relatively larger actuation travels rapidly and safely.\nYet another object of the invention is to provide a device whose electromagnetic displacement drive device provides the necessary dwell pressure and/or closing forces."}
{"text": "The present invention relates generally to clothes washing machines, and more particularly to a rotatable subassembly for a horizontal axis clothes washing machine, such assembly including a clothes basket and balance rings.\nConventional clothes washing machines include vertical axis clothes washing machines and horizontal axis clothes washing machines, where \"vertical axis\" and \"horizontal axis\" refer to the orientation of the axis of rotation of the clothes basket. The clothes basket is generally enclosed by, and rotatably attached to, the non-rotating tub, and the tub is generally enclosed by, and typically suspended from, the cabinet.\nConventional vertical axis clothes washing machines include those in which the tub is suspended from the cabinet by a rod and spring tub suspension subassembly which prevents the tub from striking the cabinet during washing. The subassembly includes a cylinder typically (pivotably) attached to the tub, a rod having a first end positioned within the cylinder and a second end (pivotably) attached to the cabinet, a piston positioned within the cylinder and attached to the first end of the rod, and a spring positioned within the cylinder between the first and second ends of the rod.\nConventional vertical axis clothes washing machines also include those having a pair of balance rings with one balance ring being attached to the top and the other balance ring being attached to the bottom of the rotating clothes basket. The balance rings have circumferential cavities which are partially filled (typically between ten and ninety percent) with water, which are sealed, and which are independent of each other. The balance rings help balance an unbalanced load of clothes during the spin-dry cycle, as is known to those skilled in the art.\nKnown horizontal axis clothes washing machines lack balance rings and typically have the tub suspended by simple springs and damped from below by friction dampers or supported from below by shock absorber type struts. Such horizontal axis clothes washing machines tend to have their cabinet feet lift and tend to \"walk\" about during the spin-dry cycle, especially with a clothes-load unbalance, unless they are weighted or bolted down to the floor. Additionally, such horizontal axis clothes washing machines can transmit significant loads to the floor during the wash cycle. What is needed is a horizontal axis clothes washing machine which encompasses all of the following characteristics during both the wash and spin-dry cycles: small excursions; small bearing loads; small dynamic load transfer to the floor; and small weight."}
{"text": "1. Field of the Invention\nThe present invention relates to an electronic-component holding apparatus, an electronic-component mounting system, and an electronic-component mounting method, and particularly to an electronic-component holding apparatus including a plural-nozzle head that holds a plurality of suction nozzles adapted to hold respective electronic components, and a head holding member that holds the plural-nozzle head such that the plural-nozzle head is detachable from the head holding member; an electronic-component mounting system including the electronic-component holding apparatus; and an electronic-component mounting method using the electronic-component holding apparatus.\n2. Discussion of Related Art\nAn example of the above-indicated electronic-component holding apparatus is disclosed by, e.g., Japanese Patent Application Publication No. 2004-221518. The electronic-component holding apparatus is employed by an electronic-component mounting system, and includes a head holding member that selectively holds an arbitrary one of different sorts of nozzle heads. The different sorts of nozzle heads hold different sorts and/or numbers of suction nozzles. Therefore, the head holding member selectively holds an appropriate one of the different nozzle heads that holds one or more suction nozzles corresponding to a sort or sorts, and/or a total number, of electronic components to be mounted on a circuit substrate.\nThe head holding member is provided with a head clamping device including an engaging pin and an engaging-pin operating device. The engaging-pin operating device includes a cam plate, and a rod member that is operable by an operator to rotate the cam plate and thereby move the engaging pin upward and downward. When a nozzle head is held by the head holding member, the operator rotates the cam plate so as to retract the engaging pin to its retracted position and, in this state, the nozzle head is attached to the head holding member. Then, the cam plate is rotated in a reverse direction so as to move the engaging pin to its operative position where the pin engages a roller provided in the nozzle head attached. When the nozzle head is removed from the head holding member, the operator operates in a reverse order, so that the engaging pin is disengaged from the roller and the nozzle head is detached from the head holding member.\nHowever, it is cumbersome to attach the nozzle head to the nozzle holding member while operating the head clamping device. That is, it is time-consuming to attach and detach, while operating the head clamping device, the nozzle head to and from the nozzle holding member, which may lead to lowering an efficiency of mounting of electronic components."}
{"text": "1. Field of the Invention\nThe present invention relates to a flat panel display (FPD) support, especially to an FPD support that allows an FPD to be stood on a plane and a visual angle of the FPD adjusted.\n2. Description of the Prior Arts\nA flat panel display (FPD) is mounted on an FPD support with a hinge. The hinge allows the FPD to pivot relative to the FPD support and to adjust a visual angle of the FPD.\nA conventional FPD support has a large base to stand stably on a plane. Consequently, the FPD with the conventional FPD support requires a large box for sale and has a large footprint, occupying a large area of a desk or table when in use.\nTo overcome the shortcomings, the present invention provides an FPD support to mitigate or obviate the aforementioned problems."}
{"text": "Microporous polymeric films having structures that enable vapors to flow through them such that they are breathable or have breathability while at the same time inhibiting or stopping the flow of liquids through them have been known for some time. These types of films have been utilized in a wide variety of applications, such as in the preparation of synthetic leather, in the preparation of cloth laminates for use as synthetic shoes, raincoats, outer wear, camping equipment such as tents, and the like, and in combination with fabrics and other materials for preparation of medical applications such as surgical gowns, bandages, and the like, and applications such as house wrap, covers for automobiles and other motorized vehicles, and the like.\nA variety of porous stretched polymeric films are known in the art. In general, there are three types of processes used to produce porous, melt-processed polymeric films including a stretching step. These three types are processes for stretching films of neat, unblended polymers containing no fillers, other than typical stabilizing additives; processes for making films from blends of two or more polymers, or from blends of polymer with mineral oil or an organic salt in which the dispersed phase is extracted with the film stretched before or after the extraction; and processes for producing films from blends of polymer with a filler such as calcium carbonate or barium sulfate with the film stretched after casting and without extraction.\nCommonly assigned U.S. Pat. No. 4,975,469 discloses oriented porous polymeric films formed from a polypropylene-based resin and having a moisture vapor transmission rate (MVTR) determined according to ASTM E-96, procedure E, in the range of about 2,500 to about 7,500 g/m.sup.2 /24 hr, and a process for forming the films. The aesthetic characteristics of these polypropylene porous films have a somewhat stiffer nature or harsher hand relative to porous film made from synthetic materials other than polypropylene and that they have a tendency to produce a \"rustling\" noise when flexed.\nCommonly assigned U.S. Pat. No. 4,386,129 discloses a porous polymeric film having pores with polygonal cross sections and an average pore diameter of about 3 to about 100 microns formed from a resinous polymer of propylene and a process for forming the film comprising the steps of forming a film containing beta-spherulites and selectively extracting the beta-spherulites. Such films have utility as filtration devices, raincoats and tents.\nJapanese Patent No. 199742 (1988) discloses a method for manufacturing a microporous polypropylene film from a composition consisting of polypropylene, a polymer with a melt crystallization point above that of polypropylene, and a beta-crystalline nucleating agent which is melt-extruded and molded into a sheet at 90.degree. to 120.degree. C., in which the resulting sheet is stretched at least uniaxially with an area draw ratio of 2.25-48.\nEP 0 273 582A discloses a process wherein polypropylene blended with mineral oil and a nucleating agent is cast into a film. The mineral oil phase separates as droplets within the polypropylene matrix, and the droplets are removed in an extraction bath. The nucleating agent is said to act to reduce the size of the droplets and thereby reduce the size of the holes in the final product. The patent also discloses that the extracted film may be biaxially stretched although no description is given regarding the moisture vapor transmission rate of the resulting material.\nU.S. Pat. No. 3,844,865 discloses a process in which polymeric materials of stretch-orientable, non-elastomeric, at least partially crystalline polymers derived from polymerization of one or more unsaturated monomers and stretch orientable polyurethanes and an inorganic salt such as calcium carbonate are formed into films and then stretched uniaxiaily or biaxially to obtain high moisture vapor transmission rates.\nU.S. Pat. No. 4,350,655 discloses a process for forming porous thermoplastic films by cold stretching at high stretch tension and low stretch ratios films formed from a blend of synthetic orientable thermoplastic polymer and at least 50 wt % of an inorganic filler such as calcium carbonate, clays and titanium oxide coated with a fatty acid ester of silicon or titanium.\nU.S. Pat. No. 4,185,148 discloses a process for producing a polypropylene film having a surface layer of beta-form crystals wherein the film is extruded with a thermal gradient such that one side of the film cools at a much more rapid rate than the other, thereby producing beta-crystals on the cooled side. The resulting film is biaxially stretched to produce a film having a rough surface and which is air-tight.\nU.S. Pat. Nos. 4,210,709, 4,331,622 and 4,335,193 disclose microporous films produced by extracting an organic substance from a film comprising 40 to 90 volume percent of a polyolefin having a number average molecular weight of 15,000 or more, 10 to 60 volume percent of an inorganic filler and 2 to 20 wt % based on the total weight of polyolefin and inorganic filler of an organic substance which is substantially insoluble in and inert to sulfuric acid and has a solubility parameter ranging from 7.3 to 8.4. Disclosed inorganic fillers are silica, calcium silicate, aluminum silicate, aluminum oxide, calcium carbonate, magnesium carbonate, kaolin clay, pulverized talc, titanium oxide and diatomaceous earth. The organic substance includes low molecular weight polyolefins such as polyethylene wax, polyisobutylene, polybutadiene, and atactic polypropylene.\nU.S. Pat. No. 4,791,144 discloses a microporous polypropylene film prepared by melt-forming into a film a mixture comprising (a) 20 to 80 wt % of a propylene homopolymer, a copolymer of propylene with other copolymerizable monomer or a blend thereof; (b) 80 to 20 wt % of at least one siliceous filler; (c) 10 to 40 wt % of a polyester type plasticizer and/or an epoxy-type plasticizer; and (d) a silane-type dispersant in an amount of 0.01 to 5 wt % based on the total weight of components (a), (b) and (c), and stretching the sheet or film at an area stretching ratio of 1.5 to 30.\nU.S. Pat. Nos. 4,777,073 and 4,929,303 disclose a breathable polyolefin film and method for preparing same from melt embossed polyolefin/filler precursor films.\nU.S. Pat. Nos. 4,814,124 and 4,921,653 disclose gas-permeable films prepared by stretching a film from a mixture containing a polyolefin resin such as polypropylene, low-density polyethylene, linear low-density polyethylene, high-density polyethylene, and polybutylene, copolymers such as ethylene-propylene copolymer and ethylene-vinyl acetate copolymer or blends of these polymers, and an inorganic filler such as calcium carbonate and barium sulfate preferably added at 50 to 500 parts by weight based on 100 parts by weight of the polyolefin resin.\nU.S. Pat. Nos. 4,824,718 and 4,902,553 disclose a method of making microporous film comprising forming a blend of a crystallizable thermoplastic polymer such as polypropylene, polyethylene and polyethylene-polypropylene, a rattle-reducing additive material and a nucleating agent system such as a solid organic acid and an inorganic compound, extruding the blend into a film and orienting the film by stretching at a temperature in the range from 10.degree. C. to a temperature 10.degree. C. below the melting temperature of the thermoplastic polymer.\nU.S. Pat. No. 4,613,643 discloses a process for forming soft porous sheets comprising 40 to 80 wt % of an inorganic filler and 60 to 20 wt % of a polyolefin-type thermoplastic elastomer composition such as ethylene-propylene copolymer, ethylene-propylene-diene copolymer, ethylene-butene copolymer and propylene-butene copolymer, and stretching the sheet at a temperature lower than the softening temperature of the polyolefin type thermoplastic elastomer.\nU.S. Pat. No. 4,833,172 discloses a method for producing microporous material comprising a matrix consisting essentially of linear ultra-high molecular weight polyolefin, a large proportion of finely divided water-soluble siliceous filler and interconnecting pores by stretching to increase the porosity and provide a region of stretch-induced molecularly oriented polymer in the matrix.\nU.S. Pat. No. 5,008,296 discloses microporous films prepared by stretching a film cast from composition of a polyolefin and high levels of CaCO.sub.3 or glass beads and calcium stearate in two directions from 1.5 to 7 times in each direction and having Gurley porosities of from 0.1 to 85 seconds.\nDespite the porous films and methods of preparing same disclosed in the patents above, there remains a need for polymeric compositions capable of being formed into microporous films with improved strength and breathability and for a facile processes for forming such microporous films. It is an object of this invention to provide improved polymeric compositions capable of being convened into oriented microporous polymeric films. Another object of this invention is to provide oriented microporous polymeric film and improved processes for preparing these films. Other objects of this invention will be apparent to persons skilled in the art from the following description and claims.\nWe have found that the objects of this invention can be attained by providing a polymeric composition comprising ethylene-propylene block copolymer, polypropylene homopolymer or random copolymer of polypropylene having up to 10 wt % of a comonomer such as ethylene or an .alpha.-olefin of 4 to 8 carbon atoms and one or more components selected from low molecular weight polypropylene, beta-spherulite nucleating agent and inorganic filler such as calcium carbonate from which oriented polymeric microporous films are produced which have a MVTR of about 500 g/m.sup.2 /24 hr or greater as determined according to ASTM E-96, procedure E. Compositions comprising about 5 to about 30 wt % of the ethylene-propylene block copolymer, about 70 to about 95 wt % of the polypropylene homopolymer or random copolymer and beta-spherulite nucleating agent are converted into oriented polymeric microporous films by a process including the extractive removal of beta-spherulites. Polymeric compositions comprising ethylene-propylene block copolymer and polypropylene homopolymer or random copolymer in a weight of about 30-95/70-5 and one or more components elected from low molecular weight polypropylene, beta-spherulite nucleating agent and calcium carbonate are convened into oriented polymeric microporous films by orienting films cast from the compositions without any extraction step.\nAdvantageously, we have found that the polymeric compositions of this invention provide oriented microporous films that exhibit improved breathability, strength, toughness and break elongation and are useful for apparel applications such as sportswear, ski clothing, lining material for sportswear and the like, medical apparel such as surgical drapes, protective gowns and the like, equipment covers, filters, housewrap, pressure sensitive labels and the like."}
{"text": "The invention is directed to optical devices and related methods. In particular, the invention provides a method and device for emitting electromagnetic radiation using nonpolar or semipolar gallium containing substrates such as GaN, AN, InN, InGaN, AlGaN, and AlInGaN, and others. The invention provides a method and device using a gallium and nitrogen containing substrate configured on the {20-21} family of planes or an off-cut of the {20-21} family of planes towards the plus or minus c-plane and/or towards the a-plane according to one or more embodiments, but there can be other configurations. Still more particularly, the invention provides a method and resulting structures that use high indium content InGaN or thick InGaN regions to facilitate manipulation of optical modes for desired optical properties in devices such as laser diodes. Such high indium content InGaN or thick InGaN regions are integrally coupled to a strain control region, which allows for such integration in an efficient and reliable manner. The invention can be applied to optical devices, lasers, light emitting diodes, solar cells, photoelectrochemical water splitting and hydrogen generation, photodetectors, integrated circuits, and transistors, among other devices.\nIn 1960, the laser was first demonstrated by Theodore H. Maiman at Hughes Research Laboratories in Malibu. This laser utilized a solid-state flash lamp-pumped synthetic ruby crystal to produce red laser light at 694 nm. By 1964, blue and green laser output was demonstrated by William Bridges at Hughes Aircraft utilizing a gas laser design called an Argon ion laser. The Ar-ion laser utilized a noble gas as the active medium and produce laser light output in the UV, blue, and green wavelengths including 351 nm, 454.6 nm, 457.9 nm, 465.8 nm, 476.5 nm, 488.0 nm, 496.5 nm, 501.7 nm, 514.5 nm, and 528.7 nm. The Ar-ion laser had the benefit of producing highly directional and focusable light with a narrow spectral output, but the wall plug efficiency was <0.1%, and the size, weight, and cost of the lasers were undesirable as well.\nAs laser technology evolved, more efficient lamp pumped solid state laser designs were developed for the red and infrared wavelengths, but these technologies remained a challenge for blue and green and blue lasers. As a result, lamp pumped solid state lasers were developed in the infrared, and the output wavelength was converted to the visible using specialty crystals with nonlinear optical properties. A green lamp pumped solid state laser had 3 stages: electricity powers lamp, lamp excites gain crystal which lases at 1064 nm, 1064 nm goes into frequency conversion crystal which converts to visible 532 nm. The resulting green and blue lasers were called “lamped pumped solid state lasers with second harmonic generation” (LPSS with SHG) had wall plug efficiency of ˜1%, and were more efficient than Ar-ion gas lasers, but were still too inefficient, large, expensive, fragile for broad deployment outside of specialty scientific and medical applications. Additionally, the gain crystal used in the solid state lasers typically had energy storage properties which made the lasers difficult to modulate at high speeds which limited its broader deployment.\nTo improve the efficiency of these visible lasers, high power diode (or semiconductor) lasers were utilized. These “diode pumped solid state lasers with SHG” (DPSS with SHG) had 3 stages: electricity powers 808 nm diode laser, 808 nm excites gain crystal, which lases at 1064 nm, 1064 nm goes into frequency conversion crystal which converts to visible 532 nm. The DPSS laser technology extended the life and improved the wall plug efficiency of the LPSS lasers to 5-10%, and further commercialization ensue into more high-end specialty industrial, medical, and scientific applications. However, the change to diode pumping increased the system cost and required precise temperature controls, leaving the laser with substantial size, power consumption while not addressing the energy storage properties which made the lasers difficult to modulate at high speeds.\nAs high power laser diodes evolved and new specialty SHG crystals were developed, it became possible to directly convert the output of the infrared diode laser to produce blue and green laser light output. These “directly doubled diode lasers” or SHG diode lasers had 2 stages: electricity powers 1064 nm semiconductor laser, 1064 nm goes into frequency conversion crystal which converts to visible 532 nm green light. These lasers designs are meant to improve the efficiency, cost and size compared to DPSS-SHG lasers, but the specialty diodes and crystals required make this challenging today. Additionally, while the diode-SHG lasers have the benefit of being directly modulate-able, they suffer from severe sensitivity to temperature which limits their application."}
{"text": "Laundry treating appliances, such as clothes washers, clothes dryers, refreshers, and non-aqueous systems, may have a configuration based on a rotating drum that defines a treating chamber in which laundry items are placed for treating. The laundry treating appliance may have a controller that implements a number of pre-programmed cycles of operation. The user typically manually selects the cycle of operation from the given pre-programmed cycles. Each pre-programmed cycle may have any number of adjustable parameters, which may be input by the user or may be set by the controller. The controller may set the parameter according to default values, predetermined values, or responsive to conditions within the treating chamber."}
{"text": "1. Field of the Invention\nThe present invention relates to a foreign matter recovering apparatus for recovering a small foreign matter present in a space inaccessible by an operator and, in particular, a space of fuel assembly for nuclear power generation.\n2. Description of the Related Art\nIt has been generally known to use, for example, a magic hand and forceps as a means for recovering foreign matter trapped in a narrow space. Further, it has also been known to use, as means for recovering a foreign matter trapped in a normally inaccessible or hard-to-access space, a medical endoscope for immediately eliminating foreign matter by a biopsy forceps as well as a vacuum for vacuum-sucking foreign matter by a nozzle from a distant location.\nThe endoscope is guided into the human body cavity to photograph the interior of the internal organs of a human being or remotely controlled from outside the human body to remove its diseased region by the forceps so that a diagnostic treatment is made from within the human body.\nIn nuclear power generation, use is made of a fuel assembly 201 as shown in FIG. 19. The fuel assembly 201 is comprised of pipe-like fuel elements 202, . . . filled with nuclear fuel and has a full length of, for example, about 4 m with clearances 203, . . . of, for example, 2 to 3 mm created among the respective fuel elements 202, . . . . FIG. 19 shows the fuel assembly comprised of 8.times.8 fuel elements.\nIn FIG. 19, reference numerals 204 and 205 show upper and lower tying plates, respectively. The fuel elements 202, . . . are held, by the upper and lower tying plates 204 and 205, at their upper and lower end portions.\nIn FIG. 19, reference numeral 206 shows a plurality of spacers (only one of them is shown) disposed in a longitudinal direction of the fuel elements 202 such that they are held partway in the longitudinal direction of the fuel elements 202, . . . . The spacer 206 includes a mechanism for individually holding the fuel elements 202, . . . .\nExternal springs 207, . . . are disposed between the upper ends of the fuel pairs 202 and the upper tying plate 204.\nThe fuel assembly 201 is held, prior to use, for example, in a fuel storage pool 208 as shown in FIG. 20, and suspended at a water depth of about 10 m.\nFurther, the fuel assembly 201 may be temporarily held in the storage pool 208 in the event of a nuclear reactor failure, etc., so that it is inspected for its defect. Therefore, the fuel assembly can be held in the pool not only before its use but also at any necessary time during its use.\nIn the case where foreign matter, such as screws and metal pieces, is trapped in the fuel assembly 201 in the pool 208, it is necessary to recover it, but the fuel assembly 201 is suspended in the storage pool at the water depth of 10 m under a high pressure environment where there exits radiation of high intensity. The clearances 203 among the fuel elements 202 are each created on the order of as small as 2 to 3 mm.\nIt is, therefore, not possible to immediately recover such foreign matter by hand. It is required that the means for recovering the foreign matter be remotely controlled at a position adequately spaced apart from the fuel assembly 201.\nIt may be considered that the aforementioned endoscope is applied to the recovery of the foreign matter in the fuel assembly 201. The medical endoscope can be directly operated, while viewing the foreign matter through an eyepiece, in which case the instrument is remotely controlled at an operation distance of, for example, about 1 m. The control of the instrument direction, driving of the forceps, etc., are done with the use of an operation force transmitted through a flexible wire. Further, the endoscope is usually so set to have 5 mm in external diameter.\nIt is, therefore, difficult for the endoscope to recover a small foreign matter at a position adequately remote from the fuel assembly 201, that is, to guide the instrument's tip into the fuel assembly 201 and set it there, guide the instrument's tip or forceps into the clearance 203 of the fuel elements 202, and accurately drive the forceps and recover the foreign matter.\nFor a means for vacuum-sucking a foreign matter without a visual inspection, it is difficult to accurately locate its vacuum nozzle at a predetermined position of an object of interest."}
{"text": "The present invention relates to a detachable biological signal measuring pad and a biological signal measuring apparatus using the same, and more particularly, to a detachable biological signal measuring pad that may greatly improve portability and ease of use, and a biological signal measuring apparatus using the same.\nWith the advent of the well-being age in recent years, concern about health is gradually growing, and especially early identification and management of the health status of individuals has been known as the most important issue to maintain a healthy life. In particular, well monitoring of biological signals from the body such as pulse or body temperature with respect to health care groups such as the elderly or the chronically ill, is to prevent misfortune such as rapid degradation of the quality of life, and therefore the importance of good understanding of one's health status has been further emphasized. However, as society is gradually aging and population of the elderly is increasing, the number of the chronically ill is increasing due to the westernization of eating habits, whereas an environment to identify and manage biological signals associated with the health of the individuals has not been significantly improved. That is, the number of persons of the health care group whose health status needs to be particularly well-managed is increasing, but an environment to frequently measure the health status of each individual is not actually established. Typically, because a person who wishes to identify their health status has to use examination equipment provided from a medical institution, it is difficult to frequently identify his/her health status, and therefore persons of the health care group having severe ups and downs in their health status have low efficiency. In recent years, home inspection equipment capable of measuring biological signals at home such as an electrocardiogram has been released, but is generally too big and heavy to move, and only measured results are obtainable, and therefore, it is difficult for a user to directly determine which action he/she should take from the measured results unless the user goes again to a medical facility to enquire and seek advice. Accordingly, there is a demand for a new type of system for measuring biological signals that may be more easily portable, and improve utilization of measured results."}
{"text": "1. Technical Field\nThe present invention relates to a piezoelectric element and a piezoelectric element-applied device that are excellent in displacement properties.\n2. Related Art\nIn general, piezoelectric elements have a piezoelectric layer with electromechanical transducing properties and two electrodes holding the piezoelectric layer therebetween. In recent years, devices (piezoelectric element-applied devices) using such a piezoelectric element as a driving source have been increasingly developed. Examples of such a piezoelectric element-applied device under development include liquid ejecting heads represented by an ink jet recording head, MEMS devices represented by a piezoelectric MEMS device, ultrasonic measuring devices represented by an ultrasonic sensor or the like, and piezoelectric actuators.\nA piezoelectric material used for such a piezoelectric layer is expected to have good piezoelectric properties, and as such a piezoelectric material, a composite oxide having a binary perovskite structure (perovskite composite oxides) containing lead titanate (PbTiO3, hereinafter, referred to as “PT”) and lead zirconate (PbZrO3, hereinafter, referred to as “PZ”), that is, lead zirconate titanate (Pb(Zr, Ti)O3, hereinafter, referred to as “PZT”) is widely known.\nIn recent years, multicomponent perovskite composite oxides containing a third component in PZT have also been developed, and as such a multicomponent PZT, PZT containing lead nickelate niobate (Pb(Ni, Nb)O3, hereinafter, referred to as “PNN”) as a third component (PNN-PZT) is known (e.g., refer to JP-A-11-129478 and JP-A-11-207957).\nPNN-PZT, however, lacks sufficient displacement properties desired for higher performance and the like of various apparatuses having a piezoelectric element-applied device mounted therein and displacement properties comparable at least to those of PZT are needed. To obtain excellent displacement properties, selection of a substitution site and a substitution element of a perovskite composite oxide, substitution amount, and the like should be further reviewed for further optimization. Such a problem is not limited to that in piezoelectric elements mounted in ink jet recording heads, and similar problems arise in elements for piezoelectric MEMS devices, ultrasonic sensors, piezoelectric actuators, and the like."}
{"text": "Wireless communication systems have been widely deployed to provide various types of communication services including voice and data services. In general, a wireless communication system is a multiple access system that supports communication among multiple users by sharing available system resources (e.g. bandwidth, transmit power, etc.) among the multiple users. The multiple access system may adopt a multiple access scheme such as Code Division Multiple Access (CDMA), Frequency Division Multiple Access (FDMA), Time Division Multiple Access (TDMA), Orthogonal Frequency Division Multiple Access (OFDMA), or Single Carrier Frequency Division Multiple Access (SC-FDMA)."}
{"text": "The present invention relates to devices for penetrating the skin to extract blood or other fluids for analysis, and more particularly, to a skin-pricking device that is capable of transporting and analyzing physiological fluid from the skin.\nMedical testing of a patient\"\"s body fluid (such as blood) is often required for diagnosis and treatment of diseases. For example, to obtain a blood sample, the current state of-art requires either venipuncture or xe2x80x9cfinger-stickxe2x80x9d using a lancet. Finger-stick tests can be done for measuring blood sugar levels by diabetics using home test kits in xe2x80x9cpoint of care diagnosis.xe2x80x9d These test kits require that a drop of blood be obtained and placed into a measurement apparatus that calculates and then displays the glucose concentration in the blood sample. To obtain the drop of blood, the user is supplied with a lancet device, which makes a skin prick (typically in the user\"\"s finger). The user then xe2x80x9cmilksxe2x80x9d the punctured spot to draw enough blood to the surface of the finger. The blood is then transferred to a test strip placed in the apparatus for measurement. In venipuncture, the blood sample is drawn in a syringe or a vial via a hypodermic needle from a vein and delivered to a measuring apparatus. In both these methods, the blood sample is injected into or applied onto a sensing device, which may be a sensor cartridge, or sensor strip (as in home glucose analysis) or a more sophisticated apparatus as used in laboratory tests. A clear disadvantage of these solutions is that they involve multiple steps, as well as other manipulations by the user. Moreover, in the finger-stick methods, the process of milking and drawing blood sample to the surface leads to changes in the analyte concentration, due to mixing of the sample with interstitial fluid or due to exposure of the sample to the ambient environment (e.g., exposure to air). For these reasons, the current methods of finger-stick sampling cannot be used to measure parameters such as blood gas contents. Due to the low efficiency of transferring the sample to the sensing apparatuses, the two step sampling and sensing methods in venipuncture typically need to draw larger blood samples than are actually required for sensing. Since pain is related to the amount of blood extracted, venipuncture inflicts significant amount of unnecessary pain and trauma to the patient. Thus, there is a need for an apparatus and method for sampling and quickly analyzing blood from a patient with minimal pain using a small volume of sample blood.\nThe present invention provides a technique for sampling and analyzing physiological fluid acquired from a tissue on the body of an animal (i.e., physiological tissue or body tissue) such as the skin (e.g., a liquid such as blood from the capillaries beneath the surface of the skin or the interstitial fluid of a tissue). In one aspect, the present invention provides a device for sampling and analyzing a fluid from the physiological tissue (e.g., from capillaries beneath the skin) of a patient by puncture. The device includes a body that includes a needle and sensors. The needle has a point for puncturing the physiological tissue and a channel in the body conducting fluid (e.g., blood) proximally from the wound into the body. Sensors in the body are accessible to blood conducted along the channel for analysis. The device can be used to sample and analyze the fluid quickly, and with minimal exposure to air. In one aspect, the present technique of sampling and analyzing a physiological fluid includes puncturing the physiological tissue with a lancing unit having a channel therein, leading through a needle to a sharp point; conducting the fluid from the physiological tissue from the sharp point up the channel; and analyzing the fluid with sensors along the channel while the fluid is in the channel.\nUsing the technique of the present invention, a user can conveniently obtain and analyze a small volume of a physiological sample (such as blood) quickly, thereby ensuring that the sample does not have exposure to air. Using a device of the present invention will enable one or more constituents (and parameters) of the patient\"\"s physiological fluid to be measured to determine the biochemical parameters (which reflect the patient\"\"s physiological condition). Such parameters may include pH, glucose concentration, blood gas content, the concentration of other blood constituents (such as , but not limited to, e.g., creatinine, potassium, chloride, calcium, blood urea nitrogen), and the like. There is no need for the user to undergo extensively invasive procedures such as arterial or venous punctures to measure such parameters, hence making self-monitoring possible for data such as blood gas contents. The convenience afforded by the present invention enables a user to more closely comply with the self-monitoring program prescribed by the health professionals. Further, the device can have a needle that is disposable and well-shielded to prevent exposure of biohazard (such as blood from a sick patient) to others. Additionally, the needle according to the present invention can be mass-produced, thus making such devices more cost-effective and available to users."}
{"text": "1) Field of the Invention\nThe field of this invention relates to devices which are for the purpose of rejuvenation of the skin of a human.\n2) Description of the Prior Art\nIt has been found in the past that a human can perform certain procedures on one's skin in order to enhance the skin to make it more healthy so that its general overall appearance is improved. For example, the application of certain cleansing creams and certain softening cremes has been known to enhance the appearance of the skin. It has also been known that applying a small electrical current to one's skin will improve the skin.\nIn the past, there has been little advancement in recent years of devices which enhance the skin with an electric charge. There has been an endless number of compositions of creams, but in actuality there has been little use of any device to apply an electric charge to one's skin. Previously, the devices which were constructed to imply an electric charge did not take into consideration as to the type of electric charge, it just being an electric charge was used. Additionally, such devices of the prior art have been constructed to be cumbersome and unattractive in appearance. It is well known that in order to entice women to use such an apparatus, with women being the normal user, the apparatus must be attractive in appearance and operate simply as possible."}
{"text": "This application claims priority to German Patent Application DE10161292.3, filed Dec. 13, 2001, the entirety of which is incorporated by reference herein.\nThis invention relates to a shroud for the support of roots of variable stator vanes in the high-pressure compressor of a gas turbine.\nMore particularly, the present invention relates to a shroud as it is used for the inner support of variable stator vanes primarily in the area of the high-pressure compressors of present-day turbofan engines.\nSuch shrouds are usually made of aluminium, steel alloy or titanium. A design of this type is exemplified in U.S. Pat. No. 5,062,767.\nThe shrouds normally consist of two parts so that there is a front and a rear shroud segment. These shroud segments are annular. For weight reduction, an axial annulus is provided in these shroud segments which is produced by turning machining, for example. In order to reduce air swirls, vibrations and noise generation, the annulus is provided with a rubber filling in the state of the art. This filling is applied by way of vulcanization.\nSuch rubber fillings are disadvantageous in several respects. One disadvantage is the high effort required for production involving a manual special process which is time-consuming and, therefore, expensive. This special process normally comprises several operations, in particular cleaning, priming, filling, heat-treatment and removal of surplus rubber material as well as subsequent re-cleaning of the segments of the shroud. Quality problems may arise from lack of adhesion or from bubbling of the rubber material. In operation, the rubber filling is susceptible to failure since it is subject to ageing and tends to develop cracks at the joints. Furthermore, the rubber filling has a relatively high weight.\nIn a broad aspect, the present invention provides a shroud of the type specified at the beginning which, while avoiding the disadvantages of the state of the art, features simplicity of design, easy and cost-effective production and, in particular, light weight.\nIt is a particular object of the present invention to provide a solution to said problems by a combination of the features described herein, with further objects and advantages of the present invention becoming apparent from the description below.\nThe present invention, therefore, provides for an annular cover arranged in the area of the opening of the annulus, this annular cover being provided with assembly openings. In addition, provision is made for a tubular support arranged in the spacing between the assembly recess of the cover and the corresponding assembly hole of the shroud segment.\nThe shroud according to the present invention features a variety of merits.\nIn accordance with the present invention, the rubber filling, which features the above-mentioned disadvantages, is replaced by an annular cover. Compared with a rubber-filled shroud segment, the shroud segment thus provided features a very similar or identical outer contour. Accordingly, the acoustic properties and the flow conditions, which may give rise to air swirls and vibrations, develop very favourably and at least achieve the values known from state-of-the-art shrouds.\nA major advantage of the solution according to the present invention is the simple and rapid manner in which it can be produced. The expensive operations required for the application of the rubber filling can now completely be dispensed with. This results in a clear cost advantage over the state of the art.\nAnother major advantage lies in the fact that no materials that are susceptible to ageing, such as rubber, are applied. This provides for increased reliability and life.\nStill another advantage is the weight reduction obtained owing to the absence of a filling with rubber or any other filling material.\nChanges to the design of the shroud segments are not required since the cover according to the present invention is fitted into the area of the mouth of the annulus. This enables existing shroud segments to be converted accordingly upon removal of the rubber filling.\nIt is also advantageous that the cover can be removed and re-installed at any time, for example for maintenance or other work.\nIn a particularly favourable form of the present invention, the cover is designed essentially as a flat ring. This ring may simply be made of sheet metal, for example aluminium or titanium. It is also particularly favourable to profile the cover at least at its rims. Such edging or flanging may give a particularly favourable clamping effect or an interference fit, thus ensuring the operational safety of the cover fitted into the annulus.\nIn order to avoid vibrations and similar occurrences, additional stiffening measures may be applied, such as beading or the like.\nThe support is preferably cup-shaped or hat-shaped. Like the cover, it can be made of sheet metal by rolling or deep-drawing or a similar shaping process. It is also particularly favourable if the free rim area of the support is profiled for retention of the cover, for example by bending or flanging. Thus, the support locates the cover and restrains it against the shroud segment, ensuring safe attachment. Restraint is applied via the fixing bolts which are also used for clamping the two shroud segments of the shroud. The supports are preferably dimensioned such that they are slightly shorter than the depth of the annulus of the shroud segments. This clearance may be 0.5 mm, for example. Accordingly, restraint is applied to the cover by way of the bolted connection.\nThe covers and the supports may be one-part or multi-part. The covers and/or supports may also be made of a plastic or similar material, for example by injection moulding, instead of a metallic material as described above."}
{"text": "1. Field of the Invention\nThe present invention relates to image equipment having image forming elements such as an image sensor element or a display element, and also to a vibrating device designed to vibrate the dust-screening member that is arranged at the front of each image forming element of such an image equipment.\n2. Description of the Related Art\nAs image equipment having image forming elements, there is known an image acquisition apparatus that has an image sensor element configured to produce a video signal corresponding to the light applied to its photoelectric conversion surface. Also known is an image projector that has a display element, such as liquid crystal element, which displays an image on a screen. In recent years, image equipment having such image forming elements have been remarkably improved in terms of image quality. If dust adheres to the surface of the image forming element such as the image sensor element or display element or to the surface of the transparent member (optical element) that is positioned in front of the image forming element, the image produced will have shadows of the dust particles. This makes a great problem.\nFor example, digital cameras of called “lens-exchangeable type” have been put to practical use, each comprising a camera body and a photographic optical system removably attached to the camera body. The lens-exchangeable digital camera is so designed that the user can use various kinds of photographic optical systems, by removing the photographic optical system from the camera body and then attaching any other desirable photographic optical system to the camera body. When the photographic optical system is removed from the camera body, the dust floating in the environment of the camera flows into the camera body, possibly adhering to the surface of the image sensor element or to the surface of the transparent member (optical element), such as a lens, cover glass or the like, that is positioned in front of the image sensor element. The camera body contains various mechanisms, such as a shutter and a diaphragm mechanism. As these mechanisms operate, they produce dust, which may adhere to the surface of the image sensor element as well.\nProjectors have been put to practical use, too, each configured to enlarge an image displayed by a display element (e.g., CRT or liquid crystal element) and project the image onto a screen so that the enlarged image may be viewed. In such a projector, too, dust may adhere to the surface of the display element or to the surface of the transparent member (optical element), such as a lens, cover glass or the like, that is positioned in front of the display element, and enlarged shadows of the dust particles may inevitably be projected to the screen.\nVarious types of mechanisms that remove dust from the surface of the image forming element or the transparent member (optical element) that is positioned in front of the image sensor element, provided in such image equipment have been developed.\nIn an electronic image acquisition apparatus disclosed in, for example, U.S. 2004/0169761 A1, a ring-shaped piezoelectric element (vibrating member) is secured to the circumferential edge of a glass plat shaped like a disc (dust-screening member). When a voltage of a prescribed frequency is applied to the piezoelectric element, the glass plat shaped like a disc undergoes a standing-wave, bending vibration having nodes at the concentric circles around the center of the glass plat shaped like a disc. This vibration removes the dust from the glass disc. The vibration (vibrational mode 1) produced by the voltage of the prescribed frequency is a standing wave having nodes at the concentric circles around the center of the disc. The dust particles at these nodes cannot be removed, because the amplitude of vibration at the nodes is small. In view of this, the glass plat shaped like a disc is vibrated at a different frequency, achieving a standing-wave vibration (vibrational mode 2) that has nodes at concentric circles different from those at which the nodes of vibrational mode 1 are located. Thus, those parts of the glass disc, where the nodes lie in vibrational mode 1, are vibrated at large amplitude.\nJpn. Pat. Appln. KOKAI Publication No. 2007-228246 discloses a rectangular dust-screening member and piezoelectric elements secured to the opposite sides of the dust-screening member, respectively. The piezoelectric elements produce vibration at a predetermined frequency, resonating the dust-screening member. Vibration is thereby achieved in such mode that nodes extend parallel to the sides of the dust-screening member. Further, as in the mechanism of U.S. 2004/0169761 A1, the dust-screening member is made to resonate at a different frequency, accomplishing a standing-wave vibrational mode, in order to change the opposition of nodes. Any one of these vibrational modes achieves bending vibration having nodes extending parallel to the sides of the dust-screening member."}
{"text": "In recent years, there is an increasing demand for highly durable concrete whose strength is highly increased, for example, by reducing a water content used in the concrete, and this tendency will further increase in the future. For reducing the water content, a polycarboxylic acid-based water reducing agent excellent in water reducing properties and flow retention is used. As the water content is reduced, however, fresh concrete viscosity (also referred to hereinafter as concrete viscosity) is increased which causes the problem of deterioration in workability and operativeness in delivery by pressurization with a pump, in placing, or in compacting into a form. This problem of increased viscosity cannot sufficiently be solved by the polycarboxylic acid-based water reducing agent, and there is a need for additives having higher effects on reduction of concrete viscosity.\nOn one hand, techniques wherein a phosphate-based additive is used as an additive for hydraulic compositions are disclosed. For example, JP-B 2508113 discloses a technique of using an additive composition containing a specific polycarboxylic acid-based cement dispersant, a specific nonionic compound, and a polyoxyalkylene alkyl ether phosphate as one compound selected from compounds used as air entraining agents. JP-A 1-219050 discloses a composition containing a phosphate of specific structure as a cement dispersant and a specific nonionic surfactant.\nJP-A 2002-121058 discloses an admixture containing a polyoxyalkylene alkyl ether phosphate and a polymer obtained by polymerizing a monomer having a specific polyoxyalkylene group."}
{"text": "The present invention relates generally to semiconductor device manufacturing and, more particularly, to an apparatus and method for improved pressure sensing control during wafer probing of integrated circuits.\nIntegrated circuits are often manufactured on a semiconductor substrate, such as a silicon wafer. Typically, a single wafer will have numerous integrated circuit (IC) devices formed thereon a lattice pattern. Each IC device has of numerous layers of circuitry and a collection of external bonding pads, which are small sites made from a conductive material such as aluminum. The bonding pads eventually serve as the device's connections to the pin leads.\nSince the packaging of a device is a costly procedure, it is desirable to test a device beforehand to avoid packaging faulty devices. The testing process involves initially establishing electrical contact between a probe card and the wafer, and thereafter running a series of tests on the devices on the wafer's surface. The probe card has a collection of individual electrical contact pins or probes that stand in for the normal pins and wire leads of a packaged device. The wafer is positioned relative to the probe card so that the contact pins on the probe card make contact with a wafer's bonding pads and probe pads while a special tester runs a series of electrical tests on the wafer's devices. A wafer prober is used to position the wafer with respect to the probe card.\nIn order to provide appropriately reliable electrical contact with a device under test (DUT), numerous and diverse types of probes have been developed in the technology, wherein it is normally recognized that each probe type necessitates the employment of a specific probing force. An insufficient contact force between the probe and the wafer will result in an unreliable electrical contact, while an excessive contact force will result in damage to the probes or contact pads on the DUT. Currently, most systems which are employed for wafer testing do not incorporate sufficiently reliable structure or testing steps which will readily or accurately determine the probe force that is being applied to each pad on the DUT during implementation of the testing procedure.\nEach probe technology has a characteristic probe compliance or spring rate, and thus the correct probe force occurs at a specific probe displacement. Consequently, a current wafer testing practice is to overdrive or displace the wafer the specified distance into the probe system. Unfortunately, the resulting forces may result in significant deflection of the probe support structure. This is especially a problem with probe arrays that incorporate a large number of probes. In this case, the amount of overdrive must be increased to overcome deflection of the support structure. The actual amount of the resultant force is not readily determinable and may be open to conjecture. Thus, a typical overdrive is ordinarily determined experimentally for one particular product, and through extrapolation or assumptions employed for all similar systems and products."}
{"text": "Some computer systems used by today's business organizations are rather complex. Such a system may include of a great number of separate components that, to some extent, must operate well with each other for the system as a whole to be functional. For example, a system may include one or more sophisticated software applications for performing business operations, such as in a customer relationship management (CRM) framework. Moreover, the system may include one or more services, which are software applications devoted to managing several business application programs, and which may be considered the basic building blocks for such business applications. In addition, many systems include one or more databases with an associated component, such as a search engine, for retrieving information from the data base. Moreover, a system may include various other components dedicated to specific tasks, such as maintaining network connections, managing user accounts, and producing graphical user interfaces.\nTypically, any of these exemplary components interacts with one or more of the other components during operation of the system. Each component may have alternative configurations, settings, and so on, that control how it operates in the system environment. It is therefore very important that integration between the various components is done properly. Also, some components can interact only with a specific version of certain software. Some or all of these configurations or settings may have default values when the system is first initiated. However, default settings may be, or become, inadequate. For example, a component that is preconfigured to work in a certain system environment may not function properly if the actual system includes components that were not taken into consideration in creating the default configuration. The need for special configuration or resetting of default values may be greater when a system includes components from several different manufacturers. A poorly-configured system may malfunction or not function as well as intended. The type of problems that occur can conceptually be categorized as either application related problems or content related problems.\nAn application related problem generally means that one or more application program does not work as intended. For example, a command or a performed operation does not produce any result at all or produces a result that is different from what it should be. As another example, something that should happen automatically in the system does not happen or occurs at the wrong time.\nA content related problem generally means that, while the system behavior may appear normal, the results are not reliable. For example, the results of a search performed on a database with known contents are not what they should be. As another example, data corruption may occur in the transfer of information between components in a system.\nThese problems also can occur after a user customizes the system. For example, adding one or more user created components to a preconfigured system may cause problems if the original components are not properly configured to work with the custom component, or by the custom component is not properly configured for the original components, or both.\nMoreover, it may be difficult to discern the exact cause of a particular problem that is observed. For example, when the system presents a collection of data that is retrieved from several different locations and it is discovered that the data is not complete, the user may not be able to readily determine whether this problem is due to a failure to retrieve data from a specific location or whether data is retrieved from the location but later lost or corrupted in transmission. To remedy an observed problem, the user may need to determine its origin and take suitable measures designed to cure it."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of tissue treatment, and more specifically to a system and method for applying reduced pressure at a tissue site.\n2. Description of Related Art\nClinical studies and practice have shown that providing a reduced pressure in proximity to a tissue site augments and accelerates the growth of new tissue at the tissue site. The applications of this phenomenon are numerous, but application of reduced pressure has been particularly successful in treating wounds. The treatment of wounds using reduced pressure is sometimes referred to in the medical community as “negative pressure tissue treatment,” “reduced pressure therapy,” or “vacuum therapy.” This type of treatment provides a number of benefits, including faster healing, and increased formulation of granulation tissue.\nReduced pressure treatment systems are often applied to large, highly exudating wounds present on patients undergoing acute or chronic care, as well as other severe wounds that are not readily susceptible to healing without application of reduced pressure. Low-severity wounds that are smaller in volume and produce less exudate have generally been treated using advanced dressings instead of reduced pressure treatment.\nCurrently, the use of reduced pressure treatment is not considered a viable or affordable option for low-severity wounds due to the manpower required to monitor and change system components, the requirement for trained medical personnel overseeing treatment, and the high cost of treatment. For example, the complexity of current reduced pressure treatment systems precludes a person with little or no specialized knowledge from administering such treatment to oneself or others. The size and power consumption characteristics of current reduced pressure treatment systems also limit the mobility of both the treatment system and the person to whom the treatment is being applied. Also, the high cost of current reduced pressure treatment systems can preclude the accessibility of such treatment systems to some users. Current reduced pressure treatment systems are also typically non-disposable after each treatment, and require electrical components or other powered devices in order to generate the reduced pressure used in treatment.\nWhile reduced pressure could be applied to low-volume and low-exudating wounds using traditional reduced pressure treatment systems, a need exists for a more simple system that allows reduced pressure treatment to be administered without specialized medical training A need further exists for a system that uses little power and is compact, allowing a user of the system to remain mobile and participate in normal day-to-day activities. Finally, a system is needed that is inexpensive so that the system can economically be used by a single patient and then disposed of following the end of treatment for that patient."}
{"text": "The invention relates to a partially spring-mounted drive, particularly for electrically-powered rail vehicles, of the type having a gear which, including an encapsulating gear housing, is seated on the wheel axle, and a pinion that is seated in the gear housing, with the coupled motor being fixedly secured to the truck and the resilient, elastic fastening of the gear housing to the truck being effected by at least two elastic spring elements that are mounted nearly vertically and diametrically opposite the pinion shaft, and whose lines of influence extend through the center of the pinion shaft. A drive of this type is known from the document DE-A-91 16 159.\nThere are numerous embodiments of partially spring-mounted drives having a transverse traction motor. In a drive according to DE-AS 12 14 263, the gear, including an encapsulating gear housing, is seated on the wheel axle, and the motor pinion or pinion is seated in the gear housing. In this instance, the motor is fixedly secured to the truck, and the gear housing is movably secured. On the one hand, a pendulum-motion suspension is effected via a connecting rod and, on the other hand, a resilient, elastic support against the truck is effected by way of two elastic spring elements that are mounted nearly vertically and diametrically (central-symmetrically) opposite the pinion shaft, and whose lines of influence extend through the center of the pinion shaft. A coupling is also typically disposed between the motor and the motor pinion or pinion.\nThe document mentioned at the outset, DE-A-91 16 159, discloses an axle drive, particularly for a rail vehicle wheel-set axle disposed in a truck, in which the axle of the drive motor extends parallel to the wheel-set axle, and the motor, which drives the wheel-set axle by way of a spur gear that is fixedly connected to the wheel-set axle, also drives a universal coupling and is fixedly connected to the truck. The wheel-set axle is resiliently connected to the truck in such a way that the spring path is larger by a multiple in the one, direction that is, the vertical direction, than in the other, direction that is, the horizontal direction. The spur gear is supported opposite the truck by means of a torque-converter bearing provided with universal joints. The axis of the torque-converter bearing intersects the axis of the drive shaft of the spur gear approximately perpendicularly, or passes close by it and extends in the direction of the largest spring path of the wheel-set axle. Thus, a fixing resilience of the universal joints at the ends of the torque-converter bearing is provided that permits, above all, a resilience in the vertical direction, while the transverse resilience permits only a very limited horizontal flexibility for the purpose of fixing the torque-converter bearing.\nThese constructions have system-stipulated disadvantages. With the increase in driving speeds and motor rpms over time in the quest for more lightweight construction, the motor bearings can be reduced in desirable fashion, but the gear bearings must retain their size due to the predetermined axle loads, or even be enlarged because of larger axle diameters. Because of high rpms in the event of a drive power outage, these bearings are in great danger because a stabilizing radial load is no longer present on account of the small gear mass that is to be provided for the high-speed range.\nIt is the object of the invention to remedy this situation and provide a simple fastening of the gear housing to the truck with which a stabilizing effect can be attained for protecting the bearings, in addition to an economical minimization of the deformation losses of the coupling."}
{"text": "This invention relates to novel glass-ceramics and, more particularly, to glass-ceramics suitable for use as a light filter and, more particularly, to glass-ceramics suitable for use as a band-pass filter and a gain flattening filter. The invention relates also to a light filter using such glass-ceramics.\nThere are light filters which cut or pass light of a specific wavelength and there are also light filters which reduce intensity of light without depending upon wavelength. The former includes a band-pass filter which passes only a specific wavelength, a notch pass filter which cuts a specific wavelength and high-pass and low-pass filters which pass only wavelengths shorter or longer than a specific wavelength. The latter includes an ND filter.\nLight filters can be classified also into an absorption type filter and an interference type filter. A representative absorption type filter is the ND filter and a representative interference type filter is the band-pass filter. A substrate made of plastic is used for absorption type filters such as those for photography. Since a substrate for light filters which are subject to a strong laser beam requires durability and heat resistance property, amorphous glass is exclusively employed for such substrate.\nThe band-pass filters are made by forming, on a substrate made of, e.g., glass, a multi-layer film of dielectric by alternately laminating an H layer of a dielectric thin film having a high refractive index and an L layer of a dielectric thin film having a low refractive index.\nIn a band-pass filter which is used for the WDM(wavelength division multiplexing) optical communication system, temperature stability of the center wavelength of the band poses a problem when a narrow band width for passing wavelengths is set for applying the band-pass filter to a wavelength of a higher density. More specifically, the band-pass filter is a sensitive element in which the center wavelength of the band varies even with a slight variation in temperature and, therefore, temperature compensation should be made by a temperature controller when the band-pass filter is used. Such temperature controller, however, cannot actually be employed because of limitation in the space where the band-pass filter is located. The temperature stability has become a matter of increasing importance since it is necessary to reduce the band width as the amount of light information increases.\nIn the past, amorphous glass has been used as a substrate for the band-pass filter. This prior art substrate is not sufficient in its compressive stress to the film and its durability since its thermal expansion property and mechanical strength are not sufficiently high. Further, amorphous glass has low mechanical strength and therefore tends to produce micro-cracks in processing with resulting cracking or chipping off of corner portions of the product which reduces the yield of the product. Moreover, in amorphous glass, a relatively large amount of alkali ingredient must be added if a high thermal expansion property is to be provided and this poses a problem of elution of alkali ingredient during and after forming of the dielectric film on the substrate. Thus, amorphous glass cannot sufficiently satisfy the demands for a substrate for a light filter, particularly a substrate for a band-pass filter.\nFor example, the glass-ceramics of a SiO2xe2x80x94Li2Oxe2x80x94MgOxe2x80x94P2O5 system disclosed in U.S. Pat. No. 5,626,935 containing lithium disilicate (Li2O.2SiO2) and xcex1-quartz (xcex1-SiO2) as main crystal phases and the size of the xcex1-quartz globular grains is controlled. No discussion, however, is made in this patent about optimum values of crystal grain diameter, degree of crystallization and mechanical strength including Young\"\"s modulus and specific gravity which are suitable for use as a light filter.\nJapanese Patent Application Laid-open Publication No. Hei 10-45426 discloses a SiO2xe2x80x94Li2Oxe2x80x94K2Oxe2x80x94MgOxe2x80x94ZnOxe2x80x94P2O5xe2x80x94Al2O3 system glass-ceramic or a SiO2xe2x80x94Li2Oxe2x80x94K2Oxe2x80x94MgOxe2x80x94ZnOxe2x80x94P2O5xe2x80x94Al2O3xe2x80x94ZrO2 system glass-ceramic suitable for laser texturing which contains, as its predominant crystal phase or phases, at least one of lithium disilicate (Li2O.2SiO2), a mixed crystal of lithium disilicate and xcex1-quartz (xcex1-SiO2), and a mixed crystal of lithium disilicate and xcex1-cristobalite (xcex1-SiO2). However, no discussion has been made in the prior art publication about optimum values of crystal grain diameter, degree of crystallization, and mechanical strength including Young\"\"s modulus and specific gravity which are suitable for use as a light filter.\nJapanese Patent Application Laid-open Publication No. Hei 9-35234 discloses a magnetic disk substrate made of a glass-ceramic of a SiO2xe2x80x94Al2O3xe2x80x94Li2O system having predominant crystal phases of lithium disilicate (Li2O.2SiO2) and xcex2-spodumene (Li2O.Al2O3.4SiO2). This glass-ceramic, however, contains, as its predominant crystal phase, xcex2-spodumene which has a negative thermal expansion characteristic which causes the substrate to have a negative thermal expansion characteristic and, in this glass-ceramic, growth of SiO2 crystals such as xcex1-quartz (xcex1-SiO2) and xcex1-cristobalite (xcex1-SiO2) which have a positive thermal expansion characteristic and thereby cause the substrate to have a positive thermal expansion characteristic is extremely restricted. It is difficult in this prior art glass-ceramic to obtain a coefficient of thermal expansion required in the present invention. Besides, since this glass-ceramic is so hard that excellent processability cannot be obtained. Further, since this glass-ceramic requires a high temperature of 820xc2x0 C. to 920xc2x0 C. for crystallization which prevents a large scale production of the product at a competitive cost. Furthermore, no discussion has been made in the publication as to crystal grain diameter, degree of crystallization and mechanical strength.\nInternational Publication WO97/01164 which includes the above described Japanese Patent Application Laid-open Publication No. Hei 9-35234 discloses a glass-ceramic for a magnetic disk in which temperature for crystallization is reduced (680xc2x0 C.-770xc2x0 C.). A sufficient improvement however cannot be achieved in this substrate. Besides, crystals grown in all examples disclosed are xcex2-eucryptite (Li2O.Al2O3.2SiO2) which has a negative thermal expansion characteristic. Further, no discussion has been made in this publication as to crystal grain diameter, degree of crystallization and mechanical strength so that this glass-ceramic has the same disadvantages as the aforementioned prior art glass-ceramics.\nIt is, therefore, an object of the invention to provide a material suitable for a substrate for a light filter which has eliminated the above described disadvantages of the prior art substrate and has a thermal expansion property which is sufficient for avoiding variation in the refractive index at a temperature at which a filter formed with a mono-layer or multilayer film is used (i.e., having a high coefficient of thermal expansion and thereby imparting compressive stress to the film to improve temperature stability of the refractive index of the film) and also has a mechanical property which imparts sufficient durability to the filter and further has excellent light transmittance.\nAccumulated studies and experiments made by the inventors of the present invention for achieving the above described objects of the invention have resulted in the finding, which has led to the present invention, that a glass-ceramic which is obtained by heat treating within a limited range a SiO2xe2x80x94Li2Oxe2x80x94K2Oxe2x80x94MgOxe2x80x94ZnOxe2x80x94P2O5xe2x80x94Al2O3xe2x80x94ZrO2 system glass contains, as its predominant crystal phase or phases, lithium disilicate (Li2O.2SiO2) or a combination of lithium disilicate and at least one crystal phase selected from the group consisting of xcex1-quartz (xcex1-SiO2), xcex1-quartz solid solution (xcex1-SiO, solid solution), xcex1-cristobalite (xcex1-SiO2) and xcex1-cristobalite solid solution (xcex1-SiO2 solid solution), has crystal grains of the predominant crystal phase or phases having a specific range of grain diameter, and has specific mechanical properties including a coefficient of thermal expansion, Young\"\"s modulus and specific gravity which are suitable for use as a light filter, more particularly, as a band-pass filter or a gain flattening filter.\nFor achieving the above described object of the invention, there is provided a glass-ceramic comprising, as a predominant crystal phase, lithium disilicate (Li2O.2SiO2) having a crystal grain diameter (average) of 0.05 xcexcm or below.\nIn another aspect of the invention, degree of crystallization of the lithium disilicate crystal phase is within a range from 3% to 20%.\nIn one aspect of the invention, the glass-ceramic has Young\"\"s modulus (GPa)/specific gravity of 37 or over.\nIn another aspect of the invention, the glass-ceramic comprises further at least one crystal phase selected from the group consisting of xcex1-quartz (xcex1-SiO2), and xcex1-quartz solid solution (xcex1-SiO2 solid solution), degree of crystallization of the xcex1-quartz (xcex1-SiO2) and the xcex1-quartz solid solution (xcex1-SiO2 solid solution) being within a range from 5% to 25% and a crystal grain diameter thereof (average) being 0.10 xcexcm or below.\nIn another aspect of the invention, the glass-ceramic comprises further at least one crystal phase selected from the group consisting of xcex1-cristobalite (xcex1-SiO2) and xcex1-cristobalite solid solution (xcex1-SiO2 solid solution), degree of crystallization of the xcex1-cristobalite (xcex1-SiO2) and the xcex1-cristobalite solid solution (xcex1-SiO2 solid solution) being within a range from 2% to 10% and a crystal grain diameter thereof (average) being 0.10 xcexcm or below.\nIn another aspect of the invention, the glass-ceramic comprises, as predominant crystal phases, lithium disilicate (Li2O.2SiO2) and at least one crystal phase selected from the group consisting of xcex1-quartz (xcex1-SiO2), xcex1-quartz solid solution (xcex1-SiO2 solid solution), xcex1-cristobalite (xcex1-SiO2) and xcex1-cristobalite solid solution (xcex1-SiO2 solid solution).\nIn another aspect of the invention, crystal grains of the predominant crystal phase or phases are fine and substantially of a globular shape.\nIn another aspect of the invention, the glass-ceramic is substantially free of Na2O and PbO.\nIn another aspect of the invention, the glass-ceramic has a surface roughness Ra (arithmetic mean roughness) after polishing of 5.0 xc3x85 or below.\nIn another aspect of the invention, the glass-ceramic has a coefficient of thermal expansion within a range from +65xc3x9710xe2x88x927/xc2x0 C. to +130xc3x9710xe2x88x927/xc2x0 C. in a temperature range from xe2x88x9250xc2x0 C. to +70xc2x0 C.\nIn another aspect of the invention, the glass-ceramic substrate has bending strength of 400 MPa or over.\nIn another aspect of the invention, the glass-ceramic has a composition which comprises in weight percent:\nIn another aspect of the invention, the glass-ceramic is obtained by heat treating a base glass for nucleation under a temperature within a range from 500xc2x0 C. to 600xc2x0 C. for one to seven hours and further heat treating the glass for crystallization under a temperature within a range from 700xc2x0 C. to 780xc2x0 C. for one to twelve hours.\nIn another aspect of the invention, there is provided a glass-ceramic substrate for a light filter using the glass-ceramic described above.\nIn still another aspect of the invention, there is provided a light filter provided by forming a multi-layer film on the glass-ceramic described above."}
{"text": "The present invention relates to an improved vibratory screening machine and feeder combination and to an oil or gas earth drilling installation utilizing said vibratory screening machine and feeder combination.\nBy way of background, a drill platform for drilling for oil or gas includes a drill rig having a drill casing and a bell nipple. Drilling mud is supplied through the drill casing for lubricating the drill and carrying drilled materials back to the surface. The drilled materials carried by the drilling mud are discharged from the bell nipple. The drilled materials may include rock, shale, clay or sand or combinations thereof. The drilling mud is separated from the drilled materials in a vibratory screening machine and reused. A conventional vibratory screening machine includes a feeder which receives the outflow from the bell nipple and an angularly adjustable vibratory frame having a screen bed on which screens are mounted. The feeder discharges the material to be screened onto the screen bed. The screen bed of the vibratory screening machine must be angularly adjusted for properly screening different compositions of drilled materials. The vibratory frame is positioned above a hopper which receives the drilling mud which passes through the screen bed. The height of the hopper determines the elevation of the screen bed, and the outflow from the feeder is above the screen bed. In the conventional screening machine the vibratory frame must have a walled construction proximate the feeder in order to retain the outflow of the feeder on the screen bed. Therefore, the outflow from the feeder has to be at least a foot above the screen bed to get over the walled structure of the vibratory frame. Also, the feeder is not physically connected to the screen bed because otherwise vibrations of the vibratory frame would be transmitted thereto. Additionally the feeder is not connected to the screen bed because the screen bed must be separate from the feeder so that it can be angularly adjusted.\nIt is conventional to conduct the outflow of the bell nipple to the screening machine by a gravity feed, especially considering that it is impractical to pump drilled materials containing rocks and shale. In order to have the desired gravity feed from the bell nipple, to the feeder of the conventional screening machine, the conventional screening machine must therefore be placed in a lowered position on the drill platform. This may be difficult to achieve because of space and structure limitations on the drill platform. Therefore, the lower the outflow from the feeder, that is, the closer the outflow from the feeder to the screen bed, the easier it will be to have the outflow of the bell nipple higher than the outflow from the feeder without requiring specialized placements of the vibratory screening machine on a drill platform."}
{"text": "The present invention relates to the bending of panels particularly panels to be used as skin for covering a cell within the structure of aircraft, rocket space vehicles or the like. The invention refers specifically to a device for bending panels of the type referred to above wherein the skin on one side carries stiffening ribs forming intersections or nodes and the device for obtaining bending and curving includes adjustable elements and deforming tool members which are operated under exertion of adequate force in order to obtain positive action under utilization of the location of adjacent nodes of the work and under consideration of changes in length wherein those ribs which run in the direction of curvature.\nPanels of the types to which the invention pertains and particularly the bending of such panels requires working certain material which has a certain elasticity and plastic behaviour. A conventional material used here is aluminum. Parts are worked in this fashion particularly are required not only to have low weight but to provide low drag resistance as far as air flow is concerned which means that the panel outside has to be as smooth as possible and must not have any bumps, corners or the like, but there should be gradual and smooth transitions and curvature wherever needed, to be as streamlined as possible in order to avoid formation of drag resistance points for air-flow. In addition, such parts should be made in an economic fashion. Panels of the type to which the invention pertain may either be provided with integrated ribs etc. or ribs, strings or the like may have been added subsequently e.g. through rivetting or the like.\nEuropean patent No. 78 891 B1 suggests a method and equipment for bending panels particularly for covering cells in aerospace vehicles wherein deformation obtains through plastic straight line length changes of the ribs extending in the direction of curving and wherein the deforming forces in the ribs are provided at such a distance from the skin itself so that the tension in the panel results in a neutral line within the skin.\nThe present invention refers specifically to improvements in methods and practice as per this particular European patent and has as its specific task to improve equipment such that in between two adjacent nodes of the stiffening ribs the panel and transverse ribs be subjected to bending. These latter ribs are subject to length extension following the deforming. After the deformation process over entire extension has been completed, these ribs are provided with a constant or uniform stress distribution."}
{"text": "1. Technical Field\nThe invention relates to speed transmission systems for automobiles, trucks, and other machines requiring the input shaft speeds to be different from the output shaft speeds, and in particular positive variable speed transmissions.\n2. Background Art\nThe art useful for the understanding, searching and examination of the invention is that of mechanisms and machine elements in general and more particularly positive variable speed changing mechanisms as classified in Class 74 subclasses 217,230,230.21, 244, 681, and 689.\n3. Disclosure of the Invention\nThe invention improves the conventional speed transmission systems by the introduction of unique moving bevelled mating plates between input and output shafts. The plates, named as \"Ruff Plates\", have unique surfaces that determines the rotational speed of one shaft to the other by the positioning of a plurality of continuous series of spaced-apart interconnected balls between the plates. The plates are thicker in the center and decrease to practically zero thickness at the peripheral edges. The mating surfaces of the plates comprises concentric series of crests and valleys increasing from the center to the peripheral edges in aligned sectors that allow the balls to pass between."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and methods for controlling a hydraulic actuator for use in an injection molding machine, and more particularly to controlling such a hydraulic actuator (both linear and rotary) with a processor which is disposed locally near the actuator and/or the hydraulic manifold.\n2. Related Art\nInjection molding machines produce great quantities of products at high speed. For example, the widely-used PET plastic drink containers are made at production rates of thousands per hour. During these high speed operations, various injection molding machine devices (such as the molding clamp assembly, the injectors, various control switches, and other machine components) are moved using a number of hydraulic actuators. Such hydraulic actuators are supplied with pressurized hydraulic fluid that causes movement of an internal diaphragm or piston which, in turn, drives the molding device. A control valve controls the flow of hydraulic fluid to the actuator to control movement of the diaphragm. FIG. 1 depicts a typical control valve. In FIG. 1, the hydraulic actuator 2 includes a piston 4 which divides a chamber into two halves 6 and 8. Movement of the piston 4 drives a load 10, which, for example, may comprise a mold and clamp mechanism. A position sensor 12 senses the position of the load 10 and provides a feedback signal to a system controller (to be discussed below).\nThe hydraulic actuator 2 has two hydraulic fluid lines 64 and 84 which allow hydraulic fluid to enter and escape from the chamber halves 6 and 8, respectively. Pressure transducers 66 and 86 respectively monitor the pressure in line 64 and 84 and provide output signals to the system controller.\nHydraulic fluid from a pressure source (typically a hydraulic fluid pump; not shown in FIG. 1) is provided to valve 14 through hydraulic line 74, while hydraulic fluid may be returned from valve 14 to hydraulic fluid storage tank (also not shown in FIG. 1) through hydraulic line 94. Pressure transducers 76 and 96 respectively monitor the pressure in lines 74 and 94 and provide output signals to the system controller.\nValve 14 controls the flow of hydraulic fluid through the chambers of hydraulic actuator 2 to move the piston 4 back and forth thus driving the load 10. Valve 14 has fluid ports A, B, P, and T which are respectively coupled to the hydraulic lines 64, 84, 74, and 94, as shown. The valve 14 has a straight flow section 142 and a cross-flow section 144 which are respectively driven by solenoids 146 and 148 in order to control the flow of fluid within the valve. For example, when the straight flow section 142 is driven to the A, B, P, and T ports, pressurized fluid will flow through lines 74 and 64 into chamber 6, driving the piston 4 toward the load 10. On the other hand, if cross-flow section 144 is driven to the ports A, B, P, and T, then pressurized hydraulic fluid will be provided through lines 74 and 84 to the chamber 8, driving the piston 4 away from the load 10.\nIn the related art, control of the hydraulic actuator 2 through the valve 14 was a relatively straightforward process. For example, U.S. Pat. No. 5,062,052 (incorporated herein by reference) discloses that such actuators may be controlled with an analog signal processor and/or a programmable logic controller which are disposed at a location remote from the injection molding actuators so that the processing circuitry is not damaged by machine heat and vibration. Typically, the analog signal processor and/or the programmable logic controller will perform closed-loop control of the actuator 2 through valve 14 in order to keep load 10 moving within the prescribed operational ranges. The analog signal processor and/or programmable logic controller will receive feedback signals from the pressure transducer units 66, 76, 86, and 96, and position information from position sensor 12 in order to control valve 14 according to a predetermined control program. The analog signal processor can also effect operational changes in the operation of the actuator 2 through command signals received through the programmable logic controller, for example to change the molding and clamping times used by load 10.\nThe programmable logic controller stores a plurality of predetermined control programs which cause the analog signal processor to control the analog devices of the injection molding machine. The programmable logic controller may also include circuitry for controlling the digital devices in the injection molding machine, for example, digital solenoid valves and proximity switches. The programmable logic controller thus controls the elements of the injection molding machine either through the analog signal processor or directly through the digital devices.\nIn the control scheme of the '052 Patent, however, the analog signal processor and the programmable logic controller are required to perform command and control operations for all of the various devices in the machine. This imposes a processing bottleneck. For example, the programmable logic controller may attempt to execute closed-loop control of multiple different analog devices at the same time. Typically, faster and more powerful processors have been used in an attempt to overcome such problems. Such expensive solutions have, nevertheless, been unable to solve the control timing problems experienced in known actuator control architectures.\nAnother problem with the known control architecture is the reliability of the analog signal processor and the programmable logic controller. If either one of these components fails, the entire machine must be stopped until a replacement is located, installed, and programmed to operate in the specified machine. Since each actuator in every machine has unique operating characteristics, the newly-installed processor(s) must be re-programed and/or re-parameterized with the operating characteristics of the corresponding actuator(s) before full-scale production can be resumed.\nFurthermore, the dedicated wiring inter-connections used for communication between the analog signal processor and the programmable logic controller to each actuator results in a plurality of wires which are difficult to install, maintain, and repair.\nThus, a need exists for a control architecture for hydraulic actuators in an injection molding machine which provides fast, flexible, and reliable control of the actuators."}
{"text": "Wireless communication networks are widely deployed to provide various communication services such as telephony, video, data, messaging, broadcasts, and so on. Such networks, which are usually multiple access networks, support communications for multiple users by sharing the available network resources. For example, one network may be a 3G (the third generation of mobile phone standards and technology), 4G, 5G, new radio (NR), or later system, which may provide network service via any one of various radio access technologies (RATs) including EVDO (Evolution-Data Optimized), 1×RTT (1 times Radio Transmission Technology, or simply 1×), W-CDMA (Wideband Code Division Multiple Access), UMTS-TDD (Universal Mobile Telecommunications System-Time Division Duplexing), HSPA (High Speed Packet Access), GPRS (General Packet Radio Service), or EDGE (Enhanced Data rates for Global Evolution). Such multiple access networks may also include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, orthogonal frequency division multiple access (OFDMA) systems, single-carrier FDMA (SC-FDMA) networks, 3rd Generation Partnership Project (3GPP) Long Term Evolution (LTE) networks, Long Term Evolution Advanced (LTE-A), and NR networks. Other examples of wireless communication networks may include WiFi (in accordance with IEEE 802.11), WiMAX (in accordance with IEEE 802.16), and Bluetooth® networks.\nA wireless communication network may include a number of base stations that can support communication for a number of mobile stations. A mobile station (MS) may communicate with a base station (BS) via a downlink and an uplink. The downlink (or forward link) refers to the communication link from the base station to the mobile station, and the uplink (or reverse link) refers to the communication link from the mobile station to the base station. A base station may transmit data and control information on the downlink to a mobile station and/or may receive data and control information on the uplink from the mobile station.\nA wireless device such as a BS or MS may include a transmitter and a receiver coupled to an antenna to support two-way communication. For data transmission, the transmitter may modulate a radio frequency (RF) carrier signal with data to obtain a modulated RF signal, amplify the modulated RF signal to obtain an output RF signal having the proper output power level, and transmit the output RF signal via the antenna to a receiving device. The transmitter may further include a directional coupler. For data reception, the receiver may obtain a received RF signal via the antenna and may condition and process the received RF signal to recover data sent by the transmitting device.\nA wireless device may include one or more transmitters and one or more receivers coupled to one or more antennas. It is desirable to implement the transmitters and receivers to achieve good performance while reducing circuitry and cost."}
{"text": "1. Field of Invention\nThe present invention relates to a cooling apparatus of an internal combustion engine, and includes not only the cooling apparatus but also an insert disposed in a water jacket of the engine and a cylinder block of the engine.\n2. Description of Related Art\nIn an internal combustion engine, a water jacket is formed around cylinder bores in a cylinder block, and engine cooling water is caused to flow in the water jacket to cool the engine which is heated due to combustion and sliding of a piston. A temperature of a cylinder bore wall is likely to be higher at an upper portion of the cylinder bore wall than at a lower portion of the cylinder bore wall. Therefore, if the cooling water is caused to flow uniformly at an upper portion and a lower portion of the water jacket for preventing the temperature of the upper portion of the cylinder bore wall from being too high, the lower portion of the cylinder bore will be over-cooled, resulting in an increase in a friction loss generated in the sliding of the piston with the cylinder bore.\nIn order to improve a temperature distribution in the cylinder bore wall, Japanese Utility Model Publication SHO 57-43338 discloses that an insert is disposed in a water jacket. By providing the insert in the water jacket, a heat removed from the cylinder bore wall can be controlled. More particularly, at a portion where the insert is provided, a flow amount is decreased and the temperature of the cylinder bore wall is maintained high (i.e., less of a cooling effect takes place). By the heat control, the temperature distribution of the cylinder bore wall is improved.\nIn order to insert the insert into the water jacket, in a cylinder block of a closed deck-type, it is conceivable to insert the insert through a water hole formed in an upper deck of the cylinder block.\nHowever, in the closed deck-type cylinder block, the water hole is small due to a structure of a core used in casting and is discontinuous in an extending direction of the water jacket. As a result, a size of the insert is also small, and a clearance between the insert and the cylinder bore wall is large. The insert also is discontinuous in the extending direction of the water jacket. Therefore, a heat amount removed from the cylinder bore wall is large, and the cylinder bore wall is still over-cooled."}
{"text": "As is well known, various tool holders have been utilized in the prior art which interface with a rotating spindle of a machine such as a milling or boring machine to securely interface a cutting tool to the machine during the cutting of a work piece. A rotary cutting machine typically includes a motor-rotated spindle to which a tool holder is attached, the tool holder being configured to accommodate a shank portion of a cutting tool which is ultimately used to cut a work piece. The attachment of the tool holder to the spindle is generally accomplished by providing a cavity in the spindle into which an upper end or shank portion of the tool holder is secured, as with an externally threaded bolt which is advanced through a portion of the spindle and is threadably received into an internally threaded bore extending axially within the shank portion of the tool holder. In most prior art tool holders, a central aperture is also formed in a lower end or mounting portion of the tool holder for receiving the shank portion of the cutting tool which is to be interfaced to the milling or other machine via the tool holder. Subsequent to the insertion of the shank portion of the cutting tool into the central aperture of the tool holder, the shank portion of the tool holder is drawn or pulled tightly into the spindle by the tightening of the above-described bolt so as to rigidly maintain the cutting tool within the tool holder.\nThough interfacing the cutting tool to the cutting machine, prior art tool holders typically possess certain deficiencies which detract from their overall utility. More particularly, while slower rotational speeds generally permit the cutting machine to perform adequately, high speed cutting, which is preferred, can cause substantial difficulty in producing a satisfactory work as a result of the development of vibratory forces that occur between the spindle and the tool holder. Specifically, at higher rotational cutting speeds, the cutting tool often begins to chatter or vibrate such that adequate tool control cannot be maintained and damage to the work piece, along with potential injury to the operator, can realistically occur. More particularly, the transfer of the harmonic resonance into the tool holder may give rise to slight movements thereof relative to the spindle, and in extreme cases may result in the tool holder loosening within the spindle. As will be recognized, the resonance of the tool holder relative to the spindle results in the cut in the work piece being substantially out of tolerance.\nA further deficiency with prior art tool holders is that the manner in which the shank portion of the cutting tool is secured within the central aperture of the tool holder often results in the non-concentric mounting of the cutting tool within the tool holder. Such non-concentric mounting is unacceptable in modem, high tolerance machining applications such as those performed on a vertical milling machine wherein minor variations in the concentricity of the cutting tool within the tool holder often cause catastrophic failure in the cutting operation.\nThe present invention addresses these and other deficiencies of the prior art tool holders by providing a tool holder which includes a dampening member for eliminating the harmonic resonance which typically occurs during the use of cutting tools in high speed milling applications. In the tool holder constructed in accordance with the present invention, the non-concentric mounting of the shank portion of the cutting tool within the tool holder is also substantially eliminated by the heat shrinking of the shank portion of the cutting tool within the tool holder. These, as well as other features and advantages attendant to the present invention will be discussed in more detail below."}
{"text": "Grills configured to cook food placed above the heat source are well known. Such a grill is disclosed in GB 2373713 and includes a series of gas jets embedded in a bed of clay litag. The heat from the burning gas rises upwards through the bed of clay litag and so cooks the food located above it. The bed acts to diffuse the heat and does not itself absorb much heat energy. However, known grills may have a bed of heat-absorbing material, for example, lava stones that are heated by the burning gas to a temperature at which the lava stones themselves emit sufficient heat radiation to contribute significantly to the cooking of the food, that is to say, this radiation, in addition to hot gases from the heat source filtering up through the lava stones, cooks the food.\nAs food cooks on the grill, material such as fat and meat juices drip onto the lava stones or clay litag and also onto the gas jets, which can easy become clogged by the falling material. The material will tend to collect in the bottom of the grill where the heat source is located. To maintain proper operation and hygiene, it is necessary to regularly clean the grill in order to avoid build-up of material and prevent the gas jets from becoming clogged. Over time, it may not be possible to remove all of the material and this may lead to uneven heating by gas jets and consequently uneven cooking of food occurs. Eventually, it may even become necessary to replace the gas jets.\nJP10038229 discloses an oven that contains sidewalls enclosing a substantially closed internal space. Within the sidewalls, a rack is provided for holding the food being cooked and, underneath the food, a number of burners are provided to heat the air inside the oven. The burners produce a flame that slants downwardly and include an overhang to prevent juices from the food being cooked falling on the burner nozzles. The flames heat the air in the oven but are not used to heat a solid incandescent radiant body.\nJP 20030 47563 discloses a closed oven for cooking food. Burners are provided in compartments adjacent to the sidewalls. The direction of the flames produced by the burners can be adjusted using air drawn in through side vents so that the flame can be made to point upwards or downwards; the object of altering the direction of the flames is to make the heating more even. The flames heat the air in the oven but are not used to heat a solid incandescent radiant body.\nJP 2003 033281 discloses a closed oven containing a number of burners for heating the interior of the oven. One set of burners directs flames into the space underneath the food being cooked and is located under a shield so that, if juices drip from the food being cooked, they do not fall on the burner. The flames heat the air in the oven but are not used to heat a solid incandescent radiant body.\nJP 2005 237465 discloses a closed oven having burners located above the food being cooked and a further set of burners extending through a side wall of the oven and directing flames to a space underneath the food being cooked. The flames heat the air in the oven but are not used to heat a solid incandescent radiant body.\nJP 2000 217718 discloses a closed oven containing burners located adjacent to the side wall of the oven. These burners produce a horizontally extending flame directed to the space underneath the food being cooked and a vertically oriented flame for heating the space above the food. The flames heat the air in the oven but are not used to heat a solid incandescent radiant body.\nU.S. Pat. No. 3,667,449 discloses a home grill or barbecue having a tray containing pumice rocks. A central burner is provided that produces flames that extend radially outwards from the burner in a horizontal direction. A domed-shaped ceramic deflector overlying the burner deflects the flames from the burner onto the pumice rocks, which will be heated so that they are incandescent. The food is located above the bed of pumice rocks and is cooked by radiant heat emitted by the incandescent pumice rocks."}
{"text": "1. Field of the Invention\nThe present invention relates to a levelling method and apparatus and, more particularly, to the improvement in a levelling method and apparatus for levelling an object of which the surface shape or surface roughness is to be measured.\n2. The Related Art\nLevelling apparatuses for controlling an object which is placed on a mounting member so as to have the desired attitude are known and they are generally used to place an object of measurement, a work, etc. on the mounting members of a surface roughness tester, a three-dimensional measuring apparatus and various working machines.\nIn the case of a surface roughness tester, for example, if the measuring plane is inclined, the base line is inclined in the result of measurement, as shown in FIG. 15 (A), which is unsuitable for obtaining the data on the surface roughness. Furthermore, if a measurable allowance must be set at a wide range in the case as shown in FIG. 15(A), the discrimination of the measurement is impaired. It is therefore necessary to adjust the mounting attitude so as to level in the direction of scanning by the probe (see FIG. 15(B)).\nIn the case of outputting the surface roughness three-dimensionally by measuring it several times by feeding the probe little by little in the direction orthogonal to the direction of scanning by the probe, the levelling in the direction orthogonal to the direction of scanning constitutes a very important precondition.\nIn a conventional levelling apparatus, however, the levelling only in one direction is automatized, and no levelling apparatus which is capable of automatic three-dimensional attitude control of an object placed on a mounting member has ever been produced.\nAs a result, the three-dimensional attitude of an object placed on the mounting member must be adjusted manually, which takes much labor and time."}
{"text": "2. Prior Art\nPrior art methods of removing a broken cork typically require that the remaining cork be pushed into the bottle. The cork is then either retrieved with various wire or net devices, or filtered from the wine.\nFIG. 1 is a side view of the side portion of a typical prior art corkscrew. In one embodiment, the current invention provides a substantially larger diameter of the spiral portion in order to avoid the crumbled middle portion of a broken cork, so that the spiral portion may be inserted near the side wall of the cork.\nFIGS. 2A-2C is a side view of various prior art cork removal devices which comprise a prior art spiral portion as illustrated in FIG. 1. FIG. 2A shows an example “screwpull” device, FIG. 2B shows an example “waiter's friend” device, and FIG. 2C shows an example “butterfly” device."}
{"text": "The invention relates to the selection of portfolios of diagnostic markers.\nA few single gene diagnostic markers such as her-2-neu are currently in use. Usually, however, diseases are not easily diagnosed with molecular diagnostics for one particular gene. Multiple markers are often required and the number of such markers that may be included in a assay based on differential gene modulation can be large, even in the hundreds of genes. It is desirable to group markers into portfolios so that the most reliable results are obtained using the smallest number of markers necessary to obtain such a result. This is particularly true in assays that contain multiple steps such as nucleic acid amplification steps."}
{"text": "1. Field of the Invention\nThe invention pertains to fire resistant firesleeve assemblies as used with hose or pipe systems, and particularly pertains to the sealing of the cut end of a firesleeve utilizing fibrous heat insulative components.\n2. Description of the Related Art\nHose systems, particularly fuel, lubricating and hydraulic hose used within aircraft engine compartments, often require fire resistant protection to minimize the likelihood of fire or high temperatures damaging the flexible elastomeric hose of the hydraulic circuits. Typically, resistance to fire and high temperature is achieved by encompassing the hose and portions of the associated hose fittings with a protective firesleeve. The firesleeve often consists of a silicone rubber tube surrounding the hose and/or fitting and the firesleeve usually includes an inner insulative material such as asbestos or glass fiber to provide insulation against the transfer of heat.\nTypically, firesleeve material is shipped and stored in indeterminate lengths, often wound upon reels, and the firesleeve is cut to the desired length to accommodate the length of hose with which it is to be used. As the cutting of the firesleeve to the desired length will expose the end of the firesleeve and the fibrous material thereof it is required that the end of the firesleeve be treated or sealed so as to discourage contamination of the firesleeve inner fibrous material and prevent the \"wicking\" of moisture, fluids, or other liquids into the fibrous material by capillary action.\nTo \"seal\" the cut ends of firesleeves, the conventional practice is to dip the cut firesleeve end into a silicone formulated \"end-dip\" liquid sealant so as to saturate the braided fibrous material adjacent the firesleeve end. Thereupon, the silicon end dip composition is dried and processed prior to the firesleeve being installed upon the associated hose. Such processing of the ends of firesleeves often requires a 24 hour delay, and the handling of a dipped firesleeve is messy, produces dripping, and such end dipping is hazardous and may not produce a liquid impervious end seal.\nFurther, the silicone formulated end dip used to seal firesleeves may be classified as a hazardous material requiring special and expensive transportation procedures. Further, end dipped firesleeves do not produce a significant frictional engagement with the hose or hose fitting unless the clamping pressures are unusually high and there is a tendency for the firesleeve to axially slide on the hose, possibly exposing a portion of the hose adjacent the fitting to fire in high temperatures.\n3. Objects of the Invention\nIt is an object of the invention to provide a fluid impervious end seal for hose or pipe firesleeve assemblies which is inexpensive, requires no dipping or application of fluids, requires no hazardous materials, and eliminates assembly and processing delays.\nAnother object of the invention is to provide a seal for the cut ends of hose or pipe firesleeves wherein the seal is formed by a flexible resilient liquid impervious cuff enclosing the firesleeve cut end which totally confines the cut end to produce an effective liquid seal, enhances fire protection, and produces an attractive aesthetically finished hose assembly."}
{"text": "1. Field of the Invention\nThis invention relates to a decoder of a convolutional coded data stream. More particularly this invention relates to an improved traceback facility for a Viterbi decoder.\n2. Description of the Related Art\nViterbi decoders are important in modern telecommunications, in which digital signals are convolutionally encoded. In a modern application, coded orthogonal frequency division multiplexing (\"COFDM\") requiring a convolutional code has been proposed in the European Telecommunications Standard DRAFT pr ETS 300 744 (May 1996), which specifies the framing structure, channel coding, and modulation for digital terrestrial television. This standard was designed to accommodate digital terrestrial television within the existing spectrum allocation for analog transmissions, yet provide adequate protection against high levels of co-channel interference and adjacent channel interference. A flexible guard interval is specified, so that the system can support diverse network configurations, while maintaining high spectral efficiency, and sufficient protection against co-channel interference and adjacent channel interference from existing Phase Alternating Line System (\"PAL\")/Sequential Coleur A'Memorie (\"SECAM\") services. Two modes of operation are defined: a \"2K mode\", suitable for single transmitter operation and for small single frequency networks with limited transmitter distances, and an \"8K mode\". The latter can be used for either single transmitter operation or for large single frequency networks. Various levels of quadrature amplitude modulation (\"QAM\") are supported, as are different inner code rates, in order to balance bit rate against ruggedness. The system is intended to accommodate a transport layer according to the Moving Picture Experts Group (\"MPEG\"), and is directly compatible with MPEG-2 coded TV signals (International Organization for Standardization/International Electrotechnical Commission (\"ISO/IEC\") Specification 13818).\nIn the noted European Telecommunications Standard data carriers a COFDM frame can be either quadrature phase shift keyed (\"QPSK\"), 16-QAM, 64-QAM, non-uniform 16-QAM, or non-uniform 64-QAM using Gray mapping.\nIn the following discussion reference may be to FIG. 1, which illustrates a simplified transition trellis diagram 167 according to a one-step Viterbi decoding process at a coding rate of 1/2, according to convolutional encoding with a constraint length K=3, and a generator polynomial EQU G(x)=(x.sup.2 +x+1, x.sup.2 +1). (1)\nThe rate 1/2 indicates that for every one bit input, the encoder generates two bits. The constraint length K is the maximum number of signals that can be used to generate an output. Using a transition trellis diagram such as diagram 167, and an incoming data sequence, it is possible to generate an output stream following a sequence of states S. In the diagram 167 a particular state S.sub.t can be represented by two bits. For example, state S.sub.t can assume the value 2 (binary 10), indicated by reference numeral 169. In the representation of diagram 167, in state S.sub.t+1 the bits of state S.sub.t are shifted by one position, and an incoming data bit occupies the rightmost (least significant bit) position. Thus, the state value 169 can legally transition to values 171 and 173 in state S.sub.t+1. For these two transitions, the convolutional encoder will produce values 175 and 177 respectively, indicated more generally as x.sub.t y.sub.t. All possible state transitions can be calculated for the encoder, i.e. given S.sub.t and data bit d.sub.t, the next state S.sub.t+1, x.sub.t and y.sub.t can be evaluated.\nThe path metric is a measurement of the likelihood that the state is on the original encoder state sequence at that time. The smaller the path metric, the more probable the state is, and vice versa. The branch metric is a measurement of the probability value attached to each branch depending on the input. The branch metric is taken as the Hamming weight, which is the number of differing bits between a received symbol xy.sub.rx and an expected symbol xy along every branch in each transition as shown in FIG. 1. Traceback is the method of going back through the trellis to find the initial state which produced the state with the smallest path metric.\nIn the preferred embodiment a 2-step decoding process is employed, corresponding to moving through the trellis two steps at a time. This doubles the time to calculate each step, and each traceback yields two bits, rather than one. However the number of calculations required at each state has also doubled, as each state now has four possible paths to be calculated. Only one path is required to be maintained in memory for each state. That path, known as the surviving path, is the one having the smallest path metric and is thus the most likely path.\nPuncturing has the effect of producing a higher rate of data transmission, as the code is more efficient. In exemplary Table 2, the convolutional encoder (not shown) encodes data to produce symbols x.sub.t and y.sub.t, which are then punctuated according to the puncturing matrix\nx: 1 0 PA1 y: 1 1\nto produce x'.sub.t and y'.sub.t, which are then retimed to be transmitted as I,Q in quadrature phase shift keyed modulation. When decoding with punctured data, the omitted bits do not contribute to the branch metric calculation.\nTABLE 1 ______________________________________ data d.sub.0 d.sub.1 d.sub.2 d.sub.3 d.sub.4 ______________________________________ xy x.sub.0 y.sub.0 x.sub.1 y.sub.1 x.sub.2 y.sub.2 x.sub.3 y.sub.3 x.sub.4 y.sub.4 x'y' x.sub.0 y.sub.0 y.sub.1 x.sub.2 y.sub.2 y.sub.3 x.sub.4 y.sub.4 IQ x.sub.0 y.sub.0 y.sub.1 x.sub.2 y.sub.2 y.sub.3 x.sub.4 y.sub.4 ______________________________________\nIn the simple example given above, branch metrics were calculated using the Hamming Weight. Significant improvements result if, instead of receiving either a 1 or 0, we receive a multiple bit representation of each x.sub.rx and y.sub.rx showing the relative likelihood of the signal being a 1 or 0. Thus, in a 16 level (4-bit) soft decoding, a 1 is represented by 15 (binary 1111).\nIn 16 level decoding, if, for example, xy.sub.rx =(3, 14) are received, the branch metrics may be calculated as shown in Table 2. When calculating new path metrics, the respective path metrics are calculated using these soft-calculated branch metrics, giving significant improvements in decoder performance. In the preferred embodiment, an 8 level (3-bit) soft decoding is used. Traceback is implemented using a systolic array, as explained below in detail.\nTABLE 2 ______________________________________ expected xy.sub.rx branch calculation result ______________________________________ 00 0 .vertline.0-3.vertline. + .vertline.0-14.vertline. 17 01 1 .vertline.0-3.vertline. + .vertline.15-14.vertline. 4 10 2 .vertline.15-3.vertline. + .vertline.0-14.vertline. 26 11 3 .vertline.15-3.vertline. + .vertline.15-14.vertline. 13 ______________________________________\nIn the preferred embodiment, data is convolutionally encoded using a constraint length K=7, which corresponds to a trellis having 64 states. A partial representation of a 2-step transition trellis diagram for this situation is illustrated in FIG. 2.\nIn the noted European Telecommunications Standard there is an outer Reed-Solomon code, and an inner punctured convolutional code, based on a mother convolutional code of rate 1/2 with 64 states, and having generator polynomials G.sub.1 =171.sub.OCT for X output and G.sub.2 =133.sub.OCT for Y output. In one mode of operation, other punctured rates of 2/3, 3/4, 5/6 and 7/8 are allowed.\nDuring the Viterbi decoding process, traceback requires a substantial amount of processing time and hardware resources. The traceback module in a Viterbi decoder also represents an important production cost."}
{"text": "Ellipsometry is a technique of non-destructive measurement allowing optical characterisation of a sample having a specular or quasi-specular surface. The term ellipsometry is used as a collective term for denoting various methods for studying the physical properties of test samples by means of measuring the amount by which they change the polarization of a polarized light beam. In principle, a beam with a known state of polarization is reflected from or transmitted through the test sample. After reflection at the test sample or after the light has passed through the test sample, the physical properties of the test sample can be ascertained from the change in polarization, using optical calculations.\nEllipsometry is based on the fact that a monochromatic electromagnetic wave changes its intensity and state of polarization if it non-perpendicularly strikes the interface between two dielectric media (for example, a substrate coated with a film). The ellipsometer linearly polarizes the light beam before it strikes the sample surface; linear (or circular) polarization indicates that the light gets only two perpendicular components having the same amplitude. The beam, after going through the interface between the surface film and substrate, is reflected and changes its polarization, i.e. both the ratio amplitude of the two components (Ψ) and their phase (Δ) are modified. The two components of the reflected light are no longer mutually perpendicular and have a different amplitude; hence the polarization becomes elliptical and the technique is called ellipsometry. The ellipsometer measures the experimental values of the two components of the reflected light, giving the values Ψ and Δ. The value Ψ is calculated by multiplying the ratio of amplitudes of the incident beam by the amplitudes of the reflected beam; the value Δ comes from the difference of their phases.\nEllipsometry can be implemented in situ and therefore is well suited to the study of the mechanisms of thin layer growth, the formation of interfaces and the control of the processes used to form these layers and interfaces. For example, during semiconductor fabrication, thin films are deposited on a wafer. The thickness of the films is critical and may be readily determined from ellipsometric measurements. There are numerous other manufacturing processes that also require accurate deposition of thin films. Ellipsometers are routinely used to measure the quality of these processes by analyzing the film properties.\nThe ellipsometric measurements can be made at one fixed wavelength or at several wavelengths (i.e. spectroscopic ellipsometry). According to the wavelength region of the source (near ultra violet, visible, near infra red, infra red, etc.) it is possible to analyze different properties of layers or of bulk materials. In the ultra violet and visible regions, the depth of penetration of the radiation is often slight. These wavelengths are often used for the study of surfaces and interfaces. Generally this does not allow access to layer volumetric properties of materials which can, however, be obtained through measurements in the infra red region. Infra red is also well suited to measurements of vibrational absorption (chemical bonds).\nIn reflection ellipsometry, the test sample has a reflective surface on which the light beam impinges with an oblique angle of incidence. The ellipsometric data obtained can be used for ascertaining the refractive index and the extinction coefficient of the surface material of the test sample. If the test sample is covered with a dielectric film, it is possible to determine the thickness and the refractive index of the film. Reflection ellipsometry is the most sensitive and most accurate method for measuring such films or layers. Instead of using a reflected beam, it is also possible to analyze the light beam which is transmitted through the test sample. This method is referred to as transmission ellipsometry. The same procedure which is used for measuring the change in polarization in transmission ellipsometry can also be employed for studying bulk properties of transparent materials, for example the birefrigence of crystals or optical rotation of a sugar solution. The latter is conventionally referred to as polarimetry, but it falls within the general term of ellipsometry. For the sake of simplicity hereinafter, reference will primary be made to reflection ellipsometry but it should be appreciated that, unless stated otherwise, the information and procedures set forth herein also apply to transmission ellipsometry and polarimetry as referred to above.\nIn ellipsometry, the physical parameters of the test sample, which are the aspects of interest, simultaneously affect the relative intensity and the phase delay of the two polarization components of the light beam. As external influences affect the two polarization components to the same degree, ellipsometry is very insensitive with respect to such external influences. That also explains the extremely high degree of accuracy of ellipsometric measuring methods in conventional laboratory equipment.\nAn ellipsometer essentially comprises a light source for emitting polarized light, a test sample and an analyser which analyses polarization of the light after it has been reflected at the test sample or after it has been transmitted through the test sample. More particularly, an ellipsometer typically comprises a light source and a detector together with two polarizers, one of which is disposed near the light source and is conventionally referred to as the polarizer, while the other is disposed in the vicinity of the analyzing means and is generally referred to as the analyzer. The test sample is disposed between the polarizer and the analyzer, and the assembly may include one or two devices for altering or modifying the polarization of the light, referred to as polarization modifying devices or polarization modulating devices. The polarization modulating devices may be polarizers which produce a partial polarization effect, birefrigent devices (which are referred to as compensators), or optical rotators and/or geometrical rotators. If, for example, the polarizer is rotated relative to the test sample, that is a geometrical rotation. The differences between the individual kinds of ellipsometers arise out of the choice of the device used for modifying the state of polarization. By virtue of the particular design construction selected, ellipsometers have different properties, for example in regard to accuracy, high measuring speed and suitability for operation with multiple wavelengths.\nBecause of the non-directional nature of optical laws, the sequence in which the optical components are disposed may be interchanged between the two polarizers. The mode of operation of the overall assembly then remains the same.\nIn order to carry out ellipsometric measurements, the surface of a sample is illuminated by a light beam and the state of polarization of an incident beam i is compared with that of the reflected beam r or the transmitted beam. A polarization electric field vector E is generally represented by its projections ES and EP, respectively perpendicular and parallel to the plane of incidence. The projections ES and EP are complex amplitudes.\nIn the domain of linear ellipsometry, the polarization reflectance ratio is traditionally defined to be the ratio of the reflection coefficients p- to s-polarized lightρ=RP/RS=(EP/ES)r/(EP/ES)i  (1)and is indicative of the modifications to the state of polarization produced by the surface being studied. ρ is a complex number, therefore two ellipsometric parameters must be obtained in order to determine ρ, which is generally represented in the form:ρ=tan Ψ·eiΔ=(EP/ES)r/(EP/ES)i  (2)\nThe two parameters Ψ and Δ describing the change in polarization are thus combined in the complex quantity ρ. The parameter Ψ is a measure of the relative intensities (the amplitude ratio) of the p-to-s polarization states of the probe light beam, and the parameter Δ is a measure of the relative phase shift between the p and s polarization states. The parameters Ψ and Δ, and hence the number ρ, depend, at the same time, on the surface properties of the sample, the angle of incidence of a beam and the measurement wavelength. The expression of Ψ and Δ, or of ρ, as a function of these parameters, is given by the equations of Fresnel, as is well known in the art.\nAs a conventional method of obtaining these ellipsometric parameters Ψ and Δ, the null ellipsometry method is known. In null ellipsometry, the change in the state of polarization which is caused by the test sample is compensated by suitable adjustment of the polarization modulating device so that the light beam is extinguished by the analyser. Adjustment to a minimum level of received intensity may be effected either manually or automatically. The measurement result is then the position of the polarization modulating device upon extinction of the light beam. In this method, a polarized beam is radiated from a light source onto a measurement target at a predetermined angle with respect to the target, and a beam reflected by the target, which is elliptically polarized, is transmitted through a λ/4 (quarter wave) plate and an analyzer to be guided to a light-receiving unit. While optical intensity signals obtained by the light-receiving unit are observed through, e.g., a measuring unit, the λ/4 plate and the analyzer are rotated to obtain a rotational angle at which the minimum optical intensity is observed. The above-mentioned ellipsometric parameters are calculated on the basis of this rotational angle, as is known in the art.\nSecond harmonic generation (SHG) and sum-frequency generation (SFG), which are the respective frequency doubling and frequency mixing of light, have emerged as powerful ellipsometric probes for characterizing oriented surface systems. By nature of the unique symmetry of SHG and SFG measurements, the detected signals are often dominated by the oriented chromophores at interfaces and are largely insensitive to greater numbers of randomly oriented species in the bulk. This symmetry condition has been exploited with great success in surface-specific spectroscopic measurements by SHG and SFG at solid/liquid, liquid/liquid, liquid/air, solid/air, semi-conductor, polymer and biological interfaces.\nThe intensity of the second order nonlinear beam generated from an oriented uniaxial organic film between two isotropic media is dependent on the nonlinear optical susceptibility (described by the rank three x(2) tensor), the incident intensity or intensities, the polarization states of the incident and exigent light, and the linear optical properties of the interfacial system. A routine approach for isolating relative values of the x(2) tensor elements is to compare the intensities of the s- or p-polarized nonlinear beams measured as functions of the polarization state(s) of the incident beam(s). For example, the expressions in Eqs. 3 and 4 below (or variations thereof) are commonly used in reflection SHG measurements of achiral uniaxial interfacial films performed using linearly polarized incident light with the plane of polarization rotated away from purely p-polarized by the angle γ.Is2ω=C·¼ sin2(2γ)|s1xXXZ|2(Iω)2  (3)Ip2ω=C·|s6xZXX+cos2γ(−s3xXXZ+s5xZXX−s6xZXX+s7xZZZ)|2(Iω)2  (4)In standard reflection and transmission measurements of achiral films far from resonance, relative values of the tensor elements are easily obtained from the ratios of the second harmonic intensities measured under different polarization conditions or from nonlinear curve-fitting.\nAs new samples, unique substrates and diverse interfacial systems are increasingly investigated by SHG and SFG, many of the simplifying assumptions routinely used to treat polarization measurements may not hold. For example, the Fresnel factors that describe the electric fields at an interface will be complex-valued in studies conducted in total internal reflection or in studies of multilayer films. The Fresnel factors will also generally be complex in systems in which the nonlinear film, the substrate, or the ambient medium absorbs either the incident or exigent light (e.g., in spectroscopic studies). Finally, if the nonlinear surface layer absorbs one or more of the incident or exigent optical frequencies, the x(2) tensor elements themselves will also be complex-valued. In any or all of these instances, even the simple expressions in Eqs. 3 and 4 (or variations thereof) describing SHG can become complicated. Expansion of the expressions in Eqs. 3 and 4 in the most general case of SHG measurements with complex fitting coefficients and complex tensor elements yields the expressions in Eqs. 5 and 6.\n                              I          s                =                ⁢                  C          ⁢                      1            4                    ⁢                                                    sin                2                            ⁡                              (                                  2                  ⁢                                                                          ⁢                  γ                                )                                      ⁡                          [                                                                                                                                                                        Re                            ⁡                                                          (                                                              s                                1                                                            )                                                                                2                                                ⁢                                                                              Re                            (                                                          χ                              XXZ                                                        )                                                    2                                                                    +                                              Im                        ⁢                                                                                                  ⁢                                                                              (                                                          s                              1 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                     )                                                            ⁢                                                              Re                                (                                                                  χ                                  ZXX                                                                )                                                                                      +                                                                                                                                                                                                          cos                              2                                                        ⁢                                                          γ                              ⁡                                                              [                                                                                                                       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             )                                                                                                                          +                                                                                                                              Re                                            (                                                                                          s                                              5                                                                                        )                                                                                    ⁢                                                                                      Im                                            (                                                                                          χ                                              ZXX                                                                                        )                 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                                  ⁡                                                                                          (                                                                                              χ                                                ZZZ                                                                                            )                                                                                                                                                                                                                                                                                                          ]                                                                                                                                                                          }                                    2                                                                                        (        6        )            \nIt is immediately apparent that extracting the real and imaginary components of the three independent tensor elements present in SHG from the second harmonic intensity measured as a function of γ is nontrivial. Even assuming the x(2) tensor elements are purely real, Eqs. 5 and 6 yield a minimum of two possible solutions for the relative tensor elements (for xZXX and xXXZ of like and opposite sign).\nAn alternative detection approach was pioneered during the emergence of nonlinear optical surface measurements of uniaxial films in which a rotating polarizer is placed between the sample and the detector. As described hereinabove, rather than compare normalized intensity ratios, the rotation angle of the polarizer that resulted in extinction of the nonlinear beam was used to determine relative values of the tensor elements and subsequently molecular orientation. Using this type of approach, the ambiguity regarding the relative sign between xZXX and xXXZ in principle can be removed. However, this rotating polarizer approach will only produce a true null provided the nonlinear beam is linearly polarized, which generally limits its reliable application to standard reflection and transmission measurements of vanishingly thin transparent films far from resonance (in which case both the x(2) tensor elements and the Fresnel factors are purely real)."}
{"text": "1. Field of the Invention\nThe present invention relates to a ferroelectric memory device employing ferroelectric capacitors, and more particularly, to a ferroelectric memory device where ROM data stored at a manufacturing process can be retained. Especially, the ferroelectric memory device according to the present invention is more suitable a non-volatile memory provided in a micro controller, for example.\n2. Description of the Related Art\nIn recent years, a ferroelectric memory device employing ferroelectric capacitors has been proposed as a non-volatile memory. The ferroelectric memory device stores and reads data by employing its hysteresis characteristic and residual polarization action included in a ferroelectric film of the ferroelectric capacitor. By applying an electric field to the ferroelectric capacitor in one direction, the ferroelectric capacitor becomes a polarizing state in one direction. Alternatively, by applying an electric field to the ferroelectric capacitor in the other direction, the ferroelectric capacitor becomes a polarizing state in the reversed direction. Such the polarizing state may be retained as a residual polarization even after an electric field applied to the ferroelectric capacitor is extinct. Therefore, the ferroelectric memory device is used as a non-volatile memory where data can be retained, even if a power is OFF.\nFIG. 9 is a circuitry diagram of a ferroelectric memory cell according to the prior art. The structure of a memory cell MC of FIG. 9 is called a 2T2C structure, in which one pair of transistors Q1, Q2 and one pair of ferroelectric capacitors C1, C2, connected to each of the transistors are included. Each of the transistors Q1, Q2 has a gate connected to a word line WL, a source or drain electrode to which one of bit line pairs BL, /BL is respectively connected. Additionally, the ferroelectric capacitors C1, C2 are connected to a plate line PL. A sense amplifier 10 is connected to a bit line pairs BL, /BL. Data is recorded by polarizing one capacitor pair C1, C2 of the memory cell MC shown in FIG. 9, and the recorded data is read out in the later-explained method.\nFIG. 10 shows a hysteresis curve of the ferroelectric film. FIG. 10 shows an applied electric field or voltage on the abscissa axis and a polarization charge on the ordinate axis. In the hysteresis curve, the polarizing state of the ferroelectric film is changed from a point K, for example, and is returned to the point K through points L, M, N.\nFIG. 11 illustrates a definition of a polarizing direction of the ferroelectric capacitor in this description. In FIG. 11, each polarizing states K, L, M, N shown in FIG. 10 are shown. A hysteresis characteristic of the ferroelectric film will be now explained according to FIGS. 10 and 11.\nAs shown in FIG. 11, when applying a downward electric field Ek to the ferroelectric capacitors C1, C2 by applying a voltage of 5V, for example, a downward polarization charge -qs of FIG. 11 is generated on the capacitors C1, C2. When the voltage applied between the capacitors C1, C2 are removed from this state K, after that, the state is moved to the state L and the polarization charge -qs remains on the capacitors C1, C2. On the other hand, when a voltage of 5V is applied for the ferroelectric capacitors C1, C2, in the upward direction of FIG. 11, the upward electric field Em is applied and the status becomes a polarizing state M of the polarization charge +qs. Even if the voltage application to the capacitor is removed from this state M, the polarizing status of polarization charge +qs can be retained on the capacitors as the state N.\nTherefore, in this description, the state K or M where the electric field or voltage is applied to the capacitors is shown by a bold arrow line, and the state L or N of residual polarization where there is no potential difference in the capacitor and the electric field is not applied is shown by a broken arrow line. The direction of arrow indicates each polarizing direction.\nFIG. 12 is a timing chart of writing and reading data to the memory cell having a 2T2C structure according to the prior art. In this timing chart, a word line WL, a plate line PL, a sense amplifier operation, a bit line pairs BL, /BL, each polarizing direction of capacitors C1, C2 are shown. FIG. 12 shows the time on the abscess axis.\nData writing and reading operation modes to the memory cell of FIG. 9 will be explained in accompanying with FIGS. 10 and 12. At first, it is assumed that data written in the ferroelectric memory cell is indefinite at time Wt0 of a write cycle. The bit line pairs BL, /BL are reset to an intermediate potential between levels H and L, and the word line WL and the plate line PL are set to L level. When the word line WL is driven to H level at time Wt1, the transistors Q1, Q2 of the memory cell become conductive, and one capacitor pair C1, C2 are respectively connected to the bit line pair BL, /BL. Then, at time Wt2, the sense amplifier 10 is activated according to the written data to set the bit lines BL and /BL respectively to H and L levels. As a result, the downward electric field is applied to the ferroelectric capacitor C1 and the state becomes the state K, which is the downward polarizing state. Then, no electric field is applied to the other ferroelectric capacitor C2 so that the polarizing direction is not changed.\nWhen the plate line PL is driven to H level at time Wt3, the capacitor C2 that is connected to the bit line /BL of L level is polarized in the other direction reversed to the capacitor C1. In other words, the capacitor C2 becomes the state M, and the capacitor C1 becomes the state L. After the plate line PL is returned to L level and the capacitor C1 is polarized again, then, the word line WL is returned to L level and the cell transistors Q1, Q2 are OFF at time Wt5. For the reason, the capacitor C becomes the downward polarizing state L, and the capacitor C2 becomes the upward polarizing state N. The polarization states are remained and are retained, even when the power is OFF.\nIn the reading operation mode, the bit line pair BL, /BL are pre-charged to 0V at time Rt0. When the word line WL is driven to H level at time Rt1 and the plate line PL is driven to H level, then, the state of the capacitor C1 is moved from the state L to the state M and its polarization is reversed. On the other hand, the state of the capacitor C2 is moved from the state N to the state M. As a result, the ferroelectric capacitor C1, of which polarization is reversed, emits more charge than that of the ferroelectric capacitor C2, of which polarization is not reversed, to each bit line, thus a predetermined potential difference is generated between the bit line pair BL, /BL.\nAt time Rt2,the plate line PL is set to L level. As a result, although the potential of the bit line pair BL, /BL is slightly pull down, the above-described potential difference can be retained. At time Rt3, the potential difference between the bit line pair BL and /BL can be detected and be amplified by activating the sense amplifier 10. As a result, the data stored in the ferroelectric capacitor can be read out through each bit line.\nSince both of the capacitors C1, C2 are in the upward polarizing state at time Rt1, the stored data is broken. Therefore, the result of amplifying the sense amplifier 10 is given to the ferroelectric capacitors C1, C2 and data is rewritten by driving the plate line PL to H and L level respectively at each time Rt4 and Rt5. When the word line WL is set to L level at time Rt6, then, a residual polarizing state according to the stored data may be retained in the capacitors of the memory cell.\nThe above-described ferroelectric memory device is used by building in a micro controller and is used as a rewritable ROM, for example. There are some cases where a program, in which a procedure of activating the micro controller is written, because the ferroelectric memory device is non-volatile. A CPU built in the micro controller executes required operations according to the program written in the ferroelectric memory.\nHowever, when using the ferroelectric memory device instead of the ROM built in the conventional micro controller, the following inconvenience will be brought.\nImmediately after finishing a manufacturing process, data in a memory becomes indefinite in a ferroelectric memory device. Therefore, data should be written to a memory in some kinds of methods. Although a write-only device can be used as this writing method, a special circuit such that the data transmitted from the write-only device is stored in the ferroelectric memory should be provided in a micro controller. Therefore, there is a demand that desired data can be pre-recorded at the manufacturing process like the conventional ROM.\nOn the other hand, once desired data is written at the manufacturing process in the conventional ROM, the data can not be rewritten after that. Therefore, when a program is recorded in a ROM of a micro controller, the program can not be changed after that. However, since the stored data can be freely changed in the ferroelectric memory device, the use of the ferroelectric memory device instead of the conventional ROM makes it possible to change the program. Conversely speaking, freely program change means that data stored at the manufacturing process may be lost, and therefore, it is also required that the lost data can be restored or recovered."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor chip package, and relates more particularly to a stacked multichip package.\n2. Description of the Related Art\nIntegration of multiple stacked chips into a single package can increase electronic component package density and reduce the signal distances between electronic components. In other words, the integration technique can not only decrease the total volume required to prepare for individual chips, but can also improve the integral performance of the chips. Among the techniques for packaging multiple stacked chips, the package technique for multiple stacked chips of the same size is one of the most common package techniques.\nIn the package technique for multiple stacked chips of the same size, a plurality of chips of the same size are stacked one upon another on a substrate and the active surfaces of the chips usually face upward for convenient to wire bond. To prevent bonding wires of a lower chip from being damaged or shorted by an upper chip, interposers are provided between the chips. The height of each interpose should be greater than the loop heights of the corresponding bonding wires. Conventionally, the interposer can be a polyimide tape, a dummy chip, a metal piece, etc.\nHowever, using interposers requires additional processes such as, for example, a chip attach adhesive coating process, a chip mounting process, and a curing process. Such additional processes increase manufacturing cost and risk lowering manufacturing yield.\nSpecifically, manufacturing yield is an important consideration in the production of multiple chip packages. When multiple chips are packaged into an encapsulation body, the result of composite manufacturing yields greatly influences manufacturing cost. In particular, when the encapsulation body includes a high-priced chip, the risk of low manufacturing yield is more important. If a low price chip in the encapsulation body malfunctions or if a chip is damaged during the chip stacking process, the high-priced chip will have to be discarded with the nonfunctional chip, and such discard causes serious impact to the manufacturing cost.\nThus, in light of the trend of multiple chip package development and the high cost risk of multiple chip packages, a low cost and high manufacturing yield multichip package structure is required by the packaging industry."}
{"text": "German publication no. DE 10 2012 206 443 A1 (Biermann et al.) discloses a planetary gear differential transmission for a motor vehicle, the planetary gear differential transmission comprising a driving gear, by means of which the torque can be introduced into the transmission; comprising a first output member and a second output member, by means of which the torque can be outputted from the transmission; comprising a differential cage, which is designed to transmit the torque from the driving gear to the first output member, wherein the differential cage is attached to the driving gear in a rotationally fixed manner; and the first output member is arranged in such a way that it can be rotated in relation to the differential cage; and comprising a rolling bearing with a bearing shell, which is attached to the differential cage in a rotationally fixed manner.\nFor example, DE 10 2011 085 119 B3 (Martini et al.), discloses differential transmissions. Thus, this document discloses a spur gear differential, comprising a planetary gear carrier and two sun gears, which can be rotated relative to the planetary gear carrier, at least one friction disk and at least one spring means, which biases the sun gears relative to the planetary gear carrier. The spring means is spring-mounted, on the one hand, against a sun gear or the planetary gear carrier and, on the other hand, against an axial bearing. The geometrical and functional features, which are disclosed in DE 10 2011 085 119 B3 (Martini et al.), shall be incorporated by reference in its entirety."}
{"text": "Many electronic devices include operating systems that can execute applications. As the popularity of applications increases, more and more developers are attracted to create applications for various platforms. Some platforms now have applications numbering in the millions."}
{"text": "An apparatus 22 for double-side polishing semiconductor wafers as shown in FIGS. 4 to 7, wherein only the main functional members are shown in FIG. 4, is well known as an apparatus for polishing the two sides of each of semiconductor wafers W at the same time.\nThe apparatus 22 shown in FIG. 4 comprises a lower polishing turn table 24 and upper polishing turn table 26 both of which are arranged in a face to face relation with each other. Lower and upper polishing pads 24a, 26a are fixedly applied on the upper face of the lower polishing turn table 24 and the lower face of the upper polishing turn table 26 respectively. The lower and upper polishing turn tables 24, 26 are respectively rotated in opposed directions by a driving means which is not shown for clarity.\nA sun gear 28 is at the center of the lower polishing turn table 24 and held on the surface thereof and an internal gear 30 in an annular shape is located near the outside of the periphery thereof. Each of wafer carriers 32 has the shape of a disc and lies between the upper surface of the lower polishing pad 24a and the lower surface of the upper polishing pad 26a. The wafer carriers 32 are moved around in a space between the lower and upper polishing pads 24a, 26a, while each of the wafer carriers 32 turns not only about its own center but revolves around the centers of the polishing pads 24a, 26a by the actions of the sun gear 28 and internal gear 30. Each of the carriers 32 has a plurality of carrier holes 34 and semiconductor wafers W to be polished are placed in the respective carrier holes 34. In polishing the wafers W, polishing slurry is supplied between a wafer W and each of the polishing pads 24a, 26a through a hole 38 communicating between the upper and lower surfaces of the upper polishing turn table 26 from a nozzle 36. As a wafer carrier 32 turns about its center and at the same time revolves around the centers of the lower and upper polishing turn tables 24, 26, the wafers W placed in the carrier holes 34 turn about the respective centers of their own and at the same time revolve around not only the centers of the polishing turn tables 24, 26 but also the center of the carrier 32, so that the wafers W are moved around between the lower and upper polishing pads 24a, 26a and the two sides of each wafer W are polished at the same time.\nThe inner peripheral edge of a carrier hole 34 has an annular member A secured along the same edge for the purpose of preventing cracking or chipping of the edge of a semiconductor wafer W. The annular member A is produced by injection molding of plastics or it may be produced from glass fiber reinforced epoxy resin composite. Structures similar to or same as the wafer carrier 32 including the annular member A are respectively disclosed in laid open publications of Japanese Utility Model Application No. 57-177641, 58-4349 and 58-59559, and laid open publication of Japanese Patent Application No. 6-226618.\nWhen double-side polishing of a semiconductor wafer W is conducted by the use of the conventional apparatus 22 for double-side polishing semiconductor wafers above mentioned, a surface of a peripheral portion of an edge of the wafer W can be polished by the function of an annular member A as a polishing buff because of a relative rotational speed of the wafer W to the carrier hole 34 due to the rotation of the wafer W within the carrier hole 34, but the upper and lower chamfered portion of the edge of the wafer W remain unpolished. Therefore, generation of particles from such portion of the edge of the wafer W cannot be prevented in the conventional double-side polishing operation. In light of the problem, it is required that the whole edge portions are polished in double-side polishing. There are additional other problems in realizing the perfect polishing of the whole edge portions in double-side polishing that the perfect edge polishing as an additional process step is a costly operation because of its technical difficulty and moreover, the main surfaces of the wafer W are contaminated during edge polishing process."}
{"text": "Monitoring littoral seas without being detected can be desirable in times of conflict. In such cases, autonomous submersible monitoring and communications systems can provide much needed intelligence. While such devices can be deployed without detection, communicating the results of monitoring by devices submerged in the sea is problematic. Sonar provides low bandwidth over short ranges and radio communications, at all but the highest powers and lowest data rates, are blocked by salt water. Effective communication requires therefore that an antenna be raised above the sea. A variety of systems have been described for raising an antenna above the sea, but they are either expensive, impractical, or readily detected, making them unsuitable for exporting information without being detected."}
{"text": "The computation of the adaptive weight vectors is one of the most difficult of signal processing algorithms used to implement large scale adaptive processing systems. This is due to extremely high computation rates, large dynamic range, and low latency requirements imposed on these systems. For typical adaptive systems, the computation of the adaptive weights requires 50 to 150 GFLOPS of computational throughput. This type of high-throughput adaptive processing calls for special purpose hardware to achieve high efficiency that fully exploits the parallel nature of the adaptive weight computation process.\nCurrent adaptive signal processing research focuses on what is known as the optimal least squares solution to the general adaptive array nulling problem. It is estimated that signal processing for the least squares problem will dominate all other functional requirements for adaptive radar systems in the near future. This dominance occurs as a direct result of the algorithm throughput requirement, which is proportional to the cube of the number of sensor elements in the system. It is estimated that, in many systems, over 90% of the total signal processing load will be dedicated to this function alone. No currently available general purpose signal processing devices or architectures will provide the throughput without the penalty of large card or module counts.\nA key application of adaptive least squares computation is in the calculation of beamforming weight vectors for optimal interference suppression (i.e., jamming, multipath) in advanced multi-channel radar, sonar, and adaptive communications systems. However, the same technology is generally applicable to any system that must quickly generate large sample covariance matrices and perform large-scale matrix inversion operations (Least Squares applications).\nPrevious techniques for finding solutions for these systems have used either a hardwired, non-programmable systolic array approach, or more commonly, programmable digital signal processor (DSP) or microprocessor approaches. The hardwired logic methods can be very efficient, but are quite inflexible. Approaches using multiple programmable DSP devices (such as those from Texas Instruments, Motorola, etc.) are very flexible, but suffer from poor efficiency and high computation latencies, primarily due to the large number of DSP devices that must be utilized to achieve the necessary arithmetic throughput. Although a processor embodying the present invention can be implemented as a highly optimized systolic array, it retains much of the flexibility of a programmable data-flow system, allowing efficient implementation of algorithm variations.\nA version of a hardwired linear systolic array of processors, called MUSE, was developed at Massachusetts Institute of Technology Lincoln Laboratory, as described in \"MUSE--A Systolic Array for Adaptive Nulling with 64 Degrees of Freedom, Using Givens Transformations and Wafer Scale Integration,\" Technical Report 886, 18, May 1990. The MUSE array specializes in real-time adaptive antenna nulling computations for a 64 degree-of-freedom (\"DOF\") sidelobe canceller radar system. The MUSE architecture consists of a linear systolic array of vector rotation cells that perform voltage domain Cholesky factorization of a sample covariance matrix using a series of Givens transformations. The MUSE then partially backsolves the Cholesky factor for an adaptive weight vector using a fixed sidelobe canceller steering constraint.\nAlthough the MUSE architecture does attack the side-lobe canceller problem, there are some limitations that prohibit its use in a more general RLS and signal processing applications. The MUSE architecture uses a fixed steering constraint, specific to a sidelobe canceller radar system, for the computation of adaptive weights. In addition, MUSE only partially solves for these adaptive weights, requiring external processing hardware to compute the final result. Further, the MUSE architecture operates only on a fixed 64 degree-of-freedom problem size, and cannot be used efficiently on other problem sizes. In order to solve general linear system problems, a method of computing weights from a general constraint is needed. If these weights could be computed within the systolic array, the dependency on external processing could be eliminated.\nAnother limitation of the MUSE architecture is its inability to perform Cholesky factor downdating. Downdating consists of subtracting off contributions of old data samples from the Cholesky factor as new ones are added. This is useful for implementing sliding window functions on the input data, a function necessary in many signal processing applications. The MUSE architecture operates in a mode where the Cholesky factor is continuously updated. A constant forgetting factor is implemented that acts as an exponentially decaying window on the input data. Although this performs adequately in some applications, it imposes limitations on the types of processing that can be performed."}
{"text": "A dynamic random access memory (DRAM) cell typically comprises a charge storage capacitor coupled to an access device such as a Metal-Oxide-Semiconductor Field Effect Transistor (MOSFET). The MOSFET functions to apply or remove charge on the capacitor, thus affecting a logical state defined by the stored charge. The amount of charge stored on the capacitor is determined by the capacitance C=εεo A/d, where ε is the dielectric constant of the capacitor dielectric, εo is the vacuum permittivity, A is the electrode (or storage node) area, and d is the interelectrode spacing. The conditions of DRAM operation, such as operating voltage, leakage rate and refresh rate, will in general mandate that a certain minimum charge be stored by the capacitor.\nIn the continuing trend to higher memory capacity, the packing density of storage cells must increase, yet each must maintain required capacitance levels. This is a crucial demand of DRAM fabrication technologies if future generations of expanded memory array devices are to be successfully manufactured. Nevertheless, in the trend to higher memory capacity, the packing density of cell capacitors has increased at the expense of available cell area. Yet, design and operational parameters determine the minimum charge required for reliable operation of the memory cell despite decreasing cell area.\nSeveral techniques have been developed to increase the capacitance of the cell capacitor without significantly affecting the cell area. For example, capacitor electrodes having textured surface morphology have been introduced to increase the interfacial area between the dielectric thin film and the adjacent electrode and, therefore, to increase the capacitance. For example, in conventional dual-sided (or double-sided) container capacitors, hemispherical grain polysilicon (HSG) has been introduced as the material of choice for the double-sided electrode because the increased surface area of the HSG electrode and of its respective interfacial area is directly proportional to the cell capacitance. However, current technologies for the formation of a double-sided HSG electrode involve using HSG on the outside of the capacitor plate. Accordingly, when two neighboring dual-sided HSG capacitors are fabricated on a DRAM memory circuit, for example, a short circuit between the two dual-sided HSG capacitors may occur and thus, may negatively affect the characteristics of the device.\nAccordingly, as memory cell density continues to increase, there is a need for an improved method for forming a dual-sided HSG container capacitor having increased capacitance per cell and low leakage, as well as a method of forming a capacitor structure that achieves high storage capacitance without increasing the size of the capacitor and without a short circuit between the capacitor structure and an additional adjacent capacitor."}
{"text": "1. Field\nThe following description relates to a network device, more particularly, to a method of searching for a signature in a payload of a packet received through a network line.\n2. Description of the Related Art\nRecently, following the ever increasing distribution of high-speed internet networks, the need for improved performance of security device performances is increasing. Conventional network devices detect intrusion and filter packets using specific signature matching in a packet. Such detection and filtering functions of the network devices can be implemented in a hardware manner, in a software manner, or in a combined hardware and software manner. In addition, a variety of requirements for signature matching increases as the network speed becomes faster and invasion methods become more and more diverse. In particular, high-speed signature matching technologies and signature matching technologies using a specific location and a specific range in a packet are increasingly required, and thus research on various methods that can satisfy such requirements has been conducted.\nIn a high-speed packet processing system such as a firewall and an intrusion detection system, a hardware-based technology such as a ternary content addressable memory (TCAM) is used to detect a match of a packet with a security policy at high speed. TCAM is used for high-speed signature searching in a payload of a packet. The high-speed signature matching is implemented by setting a signature in a TCAM and matching a header and the payload of the packet. Based on the signature matching result, a signature range searching process is performed in hardware and software manners to determine whether the matched signature is located on a specific position in the packet or within a specific range in the packet. Due to the signature range searching process after the TCAM matching, it takes a substantial amount of time for this process, thus impeding high-speed performance of a network device."}
{"text": "1. Field of the Invention\nThe present invention relates to power-on detection circuits, and more particularly to power-on detection circuits for detecting a minimum operational frequency required to power-on a device.\n2. Description of the Prior Art\nMost circuits require certain minimum power levels for operation. It is therefore necessary to have some form of power detection circuit that can detect when this power level is reached, thereby determining that power-on is possible. If the power level detected is too low then the circuit will be unable to operate. If the power level detected is too high then some potential operational time, energy, and potential for lower voltage operation will have been wasted.\nPlease refer to FIG. 1. FIG. 1 is a diagram of a conventional power detection circuit 100. The circuit 100 comprises a first resistor R1 (12) coupled to a voltage supply, a second resistor R2 (14) coupled between R1 and ground, and a transistor 16 where the gate of the transistor 16 is coupled between R1 and R2. The source of the transistor 16 is coupled to a third resistor R3 (18) and an output device 20 and the drain of the transistor 16 is coupled to ground. The circuit 100 works by utilizing a ratio of the first resistance to the second resistance (R1/R2). When the ratio is above a certain voltage Vt it will power-on the transistor 16, causing the transistor 16 to pull down against the resistor R3, thereby generating a signal in the output device 20 that signifies that power is sufficient for power-on.\nAs the circuit 100 utilizes a ratio of two resistors, the level at which it is determined that power-on is possible is not as precise as desired. It is the aim of the present invention to provide a circuit and a method of power detection that can determine a minimum power required for power-on with a high level of accuracy."}
{"text": "1. Field of the Invention\nThe present invention relates to adhesive tape dispensing technology and more particularly, to a trigger-operated adhesive tape dispenser.\n2. Description of the Related Art\nWhen using a conventional adhesive tape dispenser, it is to adhere one end of the adhesive tape of the loaded adhesive tape roll to the surface area of the object to the packed, and then move the adhesive tape dispenser to pull the adhesive tape out of the adhesive tape roll subject to the desired length. During rotation of the tae holder wheel to let off the adhesive tape of the adhesive tape roll, the adhesive tape is smoothly adhered to the surface of the object. At final, use the cutter blade to cut off the applied adhesive tape. Thus, the use of the adhesive tape dispenser can effectively simplify the packing process and enhance the packing efficiency.\nHowever, before cutting off the applied adhesive tape, the user must rotate the tape dispenser and apply a certain amount of downward pressure to the adhesive tape dispenser to force the roller against the adhesive tape and to increase the tension of the adhesive tape that is been pulled out, enabling the cutter blade to cut off the applied adhesive tape. If the downward pressure applied to adhesive tape dispenser is insufficient or the direction of the applied downward pressure is not correct, the applied adhesive tape will not be easily cut off by the cutter blade. When this problem occurs, the user must repeatedly adjust the applied pressure and angle till that the adhesive tape can be smoothly cut off, Thus, the operation of the aforesaid conventional adhesive tape dispenser is inconvenient and requires much effort."}
{"text": "In recent years, demands for a high-resolution and high-quality image have increased in various fields of applications. As an image has a higher resolution and higher quality, an amount of data on the image increases more and more. Accordingly, when video data is transferred using media such as existing wired or wireless broadband lines or video data is stored in existing storage media, the transfer cost and the storage cost of data increase.\nIn order to effectively transfer, store, and reproduce information on high-resolution and high-quality video, high-efficiency video compressing techniques can be utilized.\nIn order to enhance video compression efficiency, a method of predicting information of a current block using information of neighboring blocks of the current block without transferring the information of the current block can be used.\nInter prediction and intra prediction can be used as the prediction method.\nIn the inter prediction, pixel values of a current picture are predicted with reference to information of other pictures. In the intra prediction, pixel values of a current picture are predicted using inter-pixel relationships in the same picture.\nWhen the inter prediction is performed, information indicating a reference picture and information indicating a motion vector from the neighboring blocks in an inter prediction mode can be utilized to designate a part of another picture to be used for the prediction."}
{"text": "Zinc, aluminium and/or combinations of aluminium and zinc (Al/Zn), are widely used as surface coatings, particularly but not exclusively for steel for protection against rust and corrosion. In practice, however, the zinc or Al/Zn coatings are susceptible to white rust or black rust respectively when exposed to the atmosphere due to reactions with moisture. Such rust is detrimental to the surface and generally makes coated steel substrates unsaleable despite the fact that the overall service life of the coated steel may remain the same and the formation of rust generally interferes with finishing operations. The ability to resist such corrosion is referred to herein as wet stack performance.\nIn order to inhibit the formation of rust on coated surfaces it is generally accepted that the treatment of a surface with a chromate imparts anti corrosive properties and this type of treatment is generally referred to as chromate passivation. However chromate is highly toxic to exposed workers and due to its high toxicity, disposal of chromium residues is difficult. Further in various markets the yellow discolouration of the treated coated surfaces is considered to be an unacceptable product attribute.\nIn order to overcome the problems associated with chromate passivation, phosphate coatings have been used. However the anticorrosion properties of phosphate have been found to be far inferior to the above mentioned chromate treatment.\nU.K. Patent Application No. 2,070,073 discloses an anti corrosive treatment for preventing white rust on galvanized steel comprising applying to the surface of a galvanized steel sheet an acidic solution containing molybdic acid or a molybdate in a concentration of 10-200 g/l calculated as molybdenum and adjusted to a pH of 1 to 6 by addition of an organic or inorganic acid. However, it has been found that the corrosion resistance of aluminium and aluminium alloy surfaces treated with the above solution is inferior to the chromate treated substrates under certain conditions and the treated surfaces suffer from an undesirable degree of discolouration. Moreover modybdate treated surfaces have been shown to change from a pale yellow/blue to a strong green colour when stored for periods of time in excess of 24 hours.\nAccordingly it is an object of the present invention to provide means for avoiding and/or ameliorating at least some of the above discussed disadvantages of the prior art."}
{"text": "This invention relates to the field of ink jet printers. More particularly, this invention relates to an improved ink jet printer cartridge and a manufacturing method therefor.\nInk jet printers require an ink supply to be in fluid contact with an ink drop ejecting device which precisely ejects ink onto a recording media. The ink supply must be maintained at a relatively constant pressure to keep the ink from discharging from the ejecting device at the wrong time or in an undesired amount. There are several methods and devices for maintaining a constant negative back-pressure including capillary fibers or foam, bladder-type configurations and lung-type configurations.\nBladder and lung configurations are preferred since they may be configured to provide a smaller volume ink reservoir for the same quantity of ink contained in a foam or fiber filled reservoir. Such bladder and lung devices may consist of a rigid frame with either one or two side panels made of a thin plastic laminate material. Both configurations depend on a flexing of the side panels which is resisted by a biasing member located either between the side panels (bladder configurations) or between one of the side panels and a rigid panel attached to the frame (lung configurations). The side panels are typically heat welded or adhesively attached to a substantially rigid frame. While heat welding is an effective, less time consuming and less costly method of attaching the flexible panels to a rigid frame, improvements in heat welding techniques are needed to increase useable product yield and provide a less fragile lung or bladder construction.\nThe foregoing and other needs are provided by an improved ink cartridge for an ink jet printer and improved manufacturing method therefor. According to the invention, a method for manufacturing a pressure control device for an ink jet cartridge is provided. The method includes providing a thermoplastic, flexible web having an edge portion, a first side and a second side and providing a substantially inflexible frame member having at least one peripheral edge. The first side of the flexible web is disposed adjacent the peripheral edge of the frame, thereby defining a substantially closed cavity. The edge portion of the web and peripheral edge of the frame are heat welded to one another using a heated platen, the heated platen having a heating surface and the peripheral edge having a bonding surface. The heat welding step is preferably conducted with a structure selected from the group consisting of an angled platen for heat welding the web and peripheral edge, a conventional platen for heat welding an angled peripheral edge and an angled platen for heat welding an angled peripheral edge whereby a bond is formed between the flexible web and the bonding surface at a temperature and pressure which minimizes thinning of the flexible web on the bonding surface at least adjacent the substantially closed cavity and wherein an angle providing the angled platen or angled peripheral edge relative to a plane defined by the flexible web ranges from about 5 to about 45 degrees.\nIn another aspect the invention provides an ink jet pen for use in an ink jet printer. The pen includes an ink jet cartridge body and ink jet cartridge attached to the cartridge body. The ink jet cartridge contains ink and a pressure control device is disposed therein for maintaining a substantially constant pressure in the ink jet cartridge. The ink jet cartridge includes a substantially inflexible frame having at least one peripheral edge. At least one thermoplastic, flexible web is provided having an edge portion, a first side and a second side. The edge portion of the first side of the flexible web is heat welded to the peripheral edge of the frame thereby defining a substantially closed cavity. The edge portion and peripheral edge are heat welded to one another at a temperature and pressure sufficient to provide a reinforced web edge portion. The heat welded flexible web is preferably provided by a heat welding structure selected from the group consisting of an angled platen for heat welding the web and peripheral edge, a conventional platen for heat welding an angled peripheral edge and an angled platen for heat welding an angled peripheral edge and wherein an angle providing the angled platen or angled peripheral edge relative to a plane defined by the flexible web ranges from about 5 to about 45 degrees.\nBy heat welding the edge portion of the flexible web and peripheral edge of the frame to one another with an angled heated platen, an angled peripheral edge, or both an angled platen and angled peripheral edge, the edge portion of the web is substantially reinforced with respect to the frame. As compared to a pressure control device for an ink jet cartridge made with a standard, non-angled platen and/or peripheral edge, the pressure control device made according to the invention withstands mishandling to a greater degree thereby providing a higher yield of usable parts. The ink cartridges containing such pressure control devices are also substantially more reliable for long term use due to less thinning of the web edge portion adjacent the closed cavity during the welding or bonding process."}
{"text": "1. Field of the Invention:\nThe invention relates to an integrated circuit for detecting a received signal, which can be used, for example, in radio receivers and in particular in mobile radio receivers, as well as a circuit configuration with an integrated circuit.\nVoltage-controlled oscillators (VCO) which are tuned with a phase locked loop (PLL) have become established in radio reception technology. The oscillator frequency of the voltage-controlled oscillator is tuned using the phase locked loop with a specific step size, for example with a frequency step size of 100 KHz in frequency-modulated radio in Europe. During the transmitter search run, the entire FM band (87.5 MHz to 108 MHz) is preferably scanned in 100 KHz increments. The oscillator frequency of the voltage-controlled oscillator fOSC is typically 10.7 MHz above the input frequency to be received, that is to say between 98.2 MHz and 118.7 MHz.\nDuring the transmitter search run, as far as possible all the transmitters which are worth receiving are to be found and their input frequency or the corresponding oscillator frequency fOSC of the voltage-controlled oscillator are to be stored. If a transmitter is detected at any particular reception frequency, the phase locked loop is stopped at the a instantaneous value by the microcontroller. The corresponding values are stored in the microcontroller so that these transmitter stations which are determined can be set directly at a later time. In order to be able to detect the transmission stations, one orxe2x80x94for sake of increased detection precisionxe2x80x94a plurality of criteria are required in order to stop the voltage-controlled oscillator at the corresponding oscillator frequency. It is therefore necessary to use one, or preferably a plurality of, evaluation criteria in the receiver in order to to detect unambiguously the presence of an input signal which is worth receiving.\nIn order to increase the reliability with which a transmitter can be detected unambiguously, it is possible to use not only the field-strength signal (level of the received signal) but also the multipath signal which indicates whether the received signal is scattered by multipath reception. To do this, the field-strength signal and the multipath signal are fed to a microprocessor which digitizes these signals and evaluates them through the use of criteria which are predefined in the microprocessor. Finally, in the microprocessor it is decided whether the transmitter is one which is worth receiving. The disadvantage of such a circuit configuration is that a plurality of lines are required from the receiver module to the microprocessor and unavoidable data traffic occurs on the data bus during the transmitter search run. However, this not only unnecessarily severely loads the microprocessor but the data traffic also constitutes a permanent source of interference in the sensitive reception system. In order to keep the data traffic as low as possible, the transmitter search criterion of intermediate frequency counting is therefore frequently not included, which however has the disadvantage of less precise transmitter identification. In many cases, for a more reliable search run stop it is not sufficient to evaluate only one of the search run stop criteria including the field-strength signal, multipath signal, zero crossover of the S curve and the limited intermediate frequency transmitted from the input mixer to the intermediate frequency amplifier. It is difficult to reliably detect weakly receivable transmitters only through the use of the field-strength.\nEuropean Patent Application No. EP 0 430 469 A2 discloses a circuit for detecting a received signal for FM receivers. The detector contains an intermediate frequency detector, a field-strength comparator and a noise signal detector whose output signals are combined with one another in a logic switching element to form a signal path. A microprocessor is driven as a function of this.\nPatent Abstracts of Japan, Volume 009, No. 248 (E-347), Oct. 4, 1985 (JP-A-60-096913) discloses how to use a noise detector to detect noise on the basis of multipath signals.\nEuropean Patent Application No. EP 0 335 141 A2 discloses a configuration for detecting various reception characteristic variables including a multipath reception, the outputs of which are coupled to a logic element.\nIt is accordingly an object of the invention to provide an integrated circuit for detecting a received signal which overcomes the above-mentioned disadvantages of the heretofore-known circuits of this general type and which permits a transmitter identification as quickly and reliably as possible while avoiding interference.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, in combination with a microprocessor, an integrated circuit for detecting a received signal, including:\nan intermediate frequency detector configured to detect an intermediate frequency and supplying a first search run stop signal if the intermediate frequency is within a given range;\na field-strength comparator configured to supply a second search run stop signal if a field-strength of the received signal exceeds a field-strength setpoint value;\na multipath comparator configured to supply a third search run stop signal if a multipath signal exceeds a given multipath setpoint valve;\na logic component operatively connected to the intermediate frequency detector, the field-strength comparator, and the multipath comparator, the logic component logically combining the first, second and third search run stop signals with one another and forming a binary stop signal, the logic component having an output and providing the binary stop signal as a statically present signal at the output of the logic component, and the microprocessor receiving, as an input signal, the binary stop signal provided at the output of the logic component;\na first analog/digital converter connected upstream of the field-strength comparator, the first analog/digital converter digitizing a field-strength signal;\na first serial/parallel converter connected between the first analog/digital converter and the field-strength comparator;\na second analog/digital converter connected upstream of the multipath comparator, the second analog/digital converter digitizing the multipath signal; and\na second serial/parallel converter connected between the second analog/digital converter and the multipath comparator.\nIn other words, the integrated circuit according to the invention for detecting a received signal has an intermediate frequency detector which supplies a first search run stop signal if the intermediate frequency lies within a specific range. Furthermore, the integrated circuit has a field-strength comparator which supplies a second search run stop signal if the field-strength of the received signal exceeds a field-strength setpoint value. The invention additionally has a multipath comparator which supplies a third search run stop signal if the multipath signal exceeds a specific multipath setpoint value. Furthermore, a logic component is provided which logically combines the three search run stop signals with one another to form a binary stop signal which is statically present at the output of the logic component and is made available to a microprocessor as its input signal.\nThe data traffic between the microprocessor and the integrated circuit is reduced, which results in a severe reduction in interference.\nThe integrated circuit has the advantage that the overall space required for the radio receiver can be reduced as a result of the integration of a first analog/digital converter, which is connected upstream of the field-strength comparator and serves to digitize the field-strength signal, into the integrated circuit. Likewise, the external electrical lines, which therefore run outside the integrated circuit, can therefore be reduced, which also reduces the susceptibility to interference.\nA second analog/digital converter which is connected upstream of the multipath comparator and which serves to digitize the multipath signal also entails the aforementioned advantages.\nThe integrated circuit has the advantage that accelerated signal processing is possible in the field-strength comparator by virtue of a first serial/parallel converter which is connected between the first analog/digital converter and the field-strength comparator.\nThis advantage also applies if a second serial/parallel converter is connected between the second analog/digital converter and the multipath comparator.\nAccording to another feature of the invention, the integrated circuit can advantageously be adapted to the ambient conditions by setting the field-strength setpoint value and the multipath setpoint value to lower values in areas or conditions with weak reception than in areas or conditions with strong reception, so that rapid transmitter signal detection is still possible.\nIn an advantageous embodiment of the integrated circuit according to the invention, the field-strength setpoint value, the multipath setpoint value and the range within which the intermediate frequency is to lie are determined by the microprocessor. This permits a flexible, rapid and simple adaptation of the setpoint values to the ambient conditions without the computing capacity of the microprocessor being appreciably restricted thereby.\nAccording to another feature of the invention, a third serial/parallel converter having an output side coupled to the field-strength comparator and to the multipath comparator is provided, the third serial/parallel converter receiving the field-strength setpoint value and the multipath setpoint value in a serial manner.\nWith the objects of the invention in view there is also provided, a circuit configuration, including:\na microprocessor having a microprocessor input terminal and at least one microprocessor output terminal, the microprocessor input terminal having a single line; and\na detection circuit including:\nat least one detection circuit input terminal;\na detection circuit output terminal connected to the single line of the microprocessor input terminal;\nan intermediate frequency detector configured to detect an intermediate frequency and supplying a first search run stop signal if the intermediate frequency is within a given range;\na field-strength comparator configured to supply a second search run stop signal if a field-strength of a received signal exceeds a field-strength setpoint value;\na multipath comparator configured to supply a third search run stop signal if a multipath signal exceeds a given multipath setpoint valve;\na logic component operatively connected to the intermediate frequency detector, the field-strength comparator, and the multipath comparator, the logic component logically combining the first, second and third search run stop signals with one another and forming a binary stop signal as a statically present output signal, the logic component providing the binary stop signal to the detection circuit output terminal, the at least one microprocessor output terminal being connected to the at least one detection circuit input terminal for transmitting the field-strength setpoint value and the multipath setpoint value as a serial data word, and the given range within which the intermediate frequency is to be;\na first analog/digital converter connected upstream of the field-strength comparator, the first analog/digital converter digitizing a field-strength signal;\na first serial/parallel converter connected between the first analog/digital converter and the field-strength comparator;\na second analog/digital converter connected upstream of the multipath comparator, the second analog/digital converter digitizing the multipath signal; and\na second serial/parallel converter connected between the second analog/digital converter and the multipath comparator.\nIn other words, a circuit configuration according to the invention includes an integrated circuit and a microprocessor in which the integrated circuit has an output terminal for the binary stop signal which is connected to an input terminal of the microprocessor which includes a single line and in which the microprocessor has at least one output terminal which is connected to at least one input terminal of the integrated circuit in order to transmit over it the field-strength setpoint value and the multipath signal value as a serial data word, and the range within which the intermediate frequency is to lie.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in an integrated circuit for detecting a received signal and a corresponding circuit configuration, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawing."}
{"text": "Solid polyethylene glycols (PEGs) are used in many pharmaceutical preparations as water soluble/absorbable excipients. However, because of their relatively lower molecular weight as compared to absorbable thermoplastic polyesters, commercially available PEGs, having a molecular weight of 10,000 or less are not suitable for use in preparing molded, or extruded, biomedical devices with high mechanical integrity, including high fracture resistance and flexural strength. Similarly, the commercially available PEGs cannot be used to form medical device coatings with sufficient flexibility and resistance to fracture to be clinically acceptable. On the other hand, exceptionally high molecular analogs of PEG, namely, polyethylene oxides having a molecular weight of 100,000 Da or more, can be used as tough coating, extrusion, or molding materials. Unfortunately, if such polyethylene oxides are used as soluble/absorbable implants, they are not fully excreted since their molecular weights exceed the cut-off molecular weight of 45,000 Da for passing through the kidneys. Other applications of polyethylene oxides having molecular weights exceeding 100,000 Da include their use (1) as highly water absorbent, swellable implants; and (2) in preparing highly viscous aqueous solutions at relatively low concentration. Once again, these applications cannot be associated with the use of these polymers as an absorbable/soluble implantable device, part of an absorbable implant, or a vehicle for a parenteral pharmaceutical formulation that are expected to absorb/biodegrade to extractable by-products. This situation evoked the need for a new approach to preparing derivatives of PEG, and allied polyalkylene glycols, that can be used effectively as coating or molded articles while allowing the release of molecular species having molecular weights that can be easily excreted through the kidneys. Accordingly, this invention deals with high to ultrahigh molecular weight derivatives of solid PEGs and allied polyalkylene glycols that can be readily absorbed and excreted."}
{"text": "As an information recording medium that record various contents such as movies, DVDs (Digital Versatile Discs) and BDs (Blu-ray (registered trademark) Discs) are mostly used.\nAn image, a voice, a caption, other graphic data and the like are recorded on each of the discs in a prescribed data recording format.\nOn the other hand, a display device, such as a television set, which displays images recorded on these discs, has been improved in resolution because of increase in the number of pixels in recent years.\nFor example, the television in the past has the number of pixels represented as the number of horizontal pixels &times; the number of vertical pixels=720 &times; 480. And the television in the past can display only SD (Standard Definition) images.\nOn the other hand, the Hi-vision television which is now widespread has the number of pixels represented as the number of horizontal pixels &times; the number of vertical pixels=1920 &times; 1080. And the Hi-vision television can display HD (High Definition) images (HD contents).\nIn addition, recently, the 4K television having four times the number of pixels in the HD image, i.e., having, for example, the number of horizontal pixels &times; the number of vertical pixels=4096 &times; 2160, and capable of displaying an image according to the so-called 4K standards (4K format) has been developed and now in a spreading stage.\nThe 4K television includes a display unit having the number of horizontal pixels &times; the number of vertical pixels=4096 &times; 2160, and is capable of displaying an image (4K image (4K contents) that is four times the HD image in resolution. In addition, a display device that displays an 8K image of a high resolution is also being studied.\nAs for contents having images of these resolutions, prescriptions in international standards (ITU) are provided with respect to a color space to be used. In other words, a color space corresponding to a color that can be output is prescribed.\nAs regards the SD image and the HD image, it is prescribed to make contents according to the following color space.\nFor the SD image, BT. 601, and\nfor the HD image, BT. 709.\nIn addition, as for the 4K image, it is demanded to support two color spaces: BT. 2020 and BT. 709. In other words, it is demanded to make contents capable of outputting colors corresponding to at least one of these two color spaces.\nNote that, BT. 709 prescribes a color space having a wider color area as compared with BT. 601. In addition, BT. 2020 prescribes a wider color space as compared with BT. 709. The color space prescribed by BT. 2020 is nearly close to a color space observed by naked eyes, and output of natural colors with a higher precision is implemented.\nNote that, display of contents having these color spaces is stated in, for example, PTL 1 (JP 2010-263598A) and the like."}
{"text": "The use of implantable stents in the vasculature and other body conduits has become commonplace since first proposed by Dotter in the 1960's. These devices were required to have a small, compacted diameter for insertion into the intended body conduit and transport, typically via a catheter, to a desired site for deployment, at which site they were expanded to a larger diameter as necessary to fit interferably with the luminal surface of the body conduit. They developed into balloon expandable stents that were expanded by plastically deforming the device with an inflatable balloon on which the expandable stent was mounted in the compacted state, the balloon being attached to the distal end of the catheter and inflated via the catheter. Self-expanding stents subsequently evolved, these devices being forcibly compacted to a small diameter and restrained at that diameter by a sleeve or other means. Following delivery to a desired site for deployment, they are released from the restraint and spring open to meet the luminal surface of the body conduit. These devices are typically made from nitinol metal alloys and typically rely on the superelastic and biocompatible character of this metal. Nitinol stents that rely on the shape memory attributes of that material are also known.\nThe evolution of implantable stents included the use of a tubular covering fitted to the stent, either to the outer or the luminal surface (or both surfaces) of the stent. These covered stents have generally come to be referred to as stent-grafts. The coverings were generally of a polymeric biocompatible material such as polyethylene terephthalate (PET) or polytetrafluoroethylene (PTFE). See, for example, U.S. Pat. No. 4,776,337 to Palmaz. This patent also describes that the covering may be optionally provided with perforations if desired for particular applications. Because of the open area provided by the perforations, such devices having perforated coverings may be considered to be a sort of hybrid stent and stent-graft, as are devices that include stent frame having metallic stent elements and polymeric elements connecting, covering or other otherwise being attached to the stent elements. The presence of the polymeric elements reduces the otherwise open space between the adjacent metallic stent elements, either very slightly or very substantially depending on the intended application and mechanical design.\nGenerally, a fully covered stent-graft can be considered to have a surface area (hereinafter Amax) equal to the circumference of the expanded stent multiplied by the length of the stent. For a conventional, open frame stent (as opposed to a stent-graft), the surface area represented by all of the stent elements is only a small portion of the maximum surface area Amax. The actual surface area covered by the stent, meaning the area covered by all components of the stent (including connecting elements) in their deployed state, is Astent. The porosity index, or P.I., describes the open area (the portion of the maximum surface area not covered by all components of the stent assembly) as a percentage of maximum surface area, wherein:P.I.=(1−(Astent/Amax))×100%.\nThe open area may be a continuous single space, such as the space between windings of a single helically wound stent element. Likewise the open area may be represented by the space between multiple individual annular or ring-shaped stent elements. The open area may also be represented by the total area of multiple apertures provided by either a single stent element (e.g., as shown by FIGS. 1B and 2B of U.S. Pat. No. 4,776,337) or by multiple stent elements providing multiple apertures. If multiple apertures are provided they may be of equal or unequal sizes. The use of a perforated graft covering or of polymeric elements in addition to metallic stent elements may also reduce the open area.\nStents having a porosity index of greater than 50% are considered to be substantially open stents.\nIn addition to the porosity index, the size of any aperture providing the open area must be considered if it is intended to cover only a portion of a stent area for a specific stent application. For multiple apertures, often the consideration must be for the largest size of any individual aperture, particularly if the apertures are to provide for a “filtering” effect whereby they control or limit the passage of biologic materials from the luminal wall into the flow space of the body conduit.\nVarious stent devices combining metallic stent elements with polymeric connecting elements are known; see, for example U.S. Pat. No. 5,507,767 to Maeda et al. Another is a stent provided with a flexible knitted sleeve having small open apertures in the fashion of chain link fencing, from InspireMD Ltd. (4 Derech Hashalom St., Tel Aviv 67892 Israel). Perforated stent-grafts are also known; see, for example WO00/42949."}
{"text": "Conventionally, forming of outer panel components of an automobile and the like have been carried out by press forming, in which a work is press processed using a press die. Generally, it is preferred that a press die used in a press process be such that the die surface, which is a contact surface with a work in the press die, be smoothed. By smoothing the die surface, it is possible to increase surface precision of a work (outer panel components) formed.\nAs a method for smoothing a die surface, the applicant proposes a processing method of a die surface and a die for which the processing method has been applied (see Patent Document 1) in which, first, a forming surface of a die is formed by a cutting process, then a synthetic resin is applied to the formed die surface, and next shot blasting is performed for smoothing a shape of the die surface. According to a method described in Patent Document 1, by a cutting process, it is possible effectively to grind high spots that are leftover from cutting generated on a die surface, thereby efficiently carrying out smoothing of a die surface.\nPatent Document 1: Japanese Examined Patent Application Publication No. Hei 8-263"}
{"text": "1. Field of the Invention\nThe present invention relates to a system for collecting scraps produced when various materials such as metals and plastics are processed by processing machines such as pressing machines and lathes and, more particularly, relates to a system which can efficiently collect scraps by the use of a suction force generated by a flow of high-pressure air.\n2. Statement of the Prior Art\nHeretofore, two typical systems have been used to collect scraps produced from a pressing machine, i.e. a press, which is one example of a machine for processing materials, when punching is carried out thereby.\nIn the first system, scraps are collected by simply receiving them as they fall due to their own weight. In this case, a box for collecting scraps is placed below the punching tool of the pressing machine, and scraps produced by punching of sheets fall due to their own weight and are received in the box and collected.\nThe second system is a dust collector. One example of such a dust collector system is disclosed in, e.g., Japanese Utility Model Laid-Open No. 61-24157. This system is designed as a dust collector for a portable air grinder, and is characterized by having an air ejector fixedly provided on a lid member of a cylindrical tank, an air filter provided in the tank at the downstream end of a suction pipe leading to the air ejector, a compressed-air inlet passage in the ejector connected with an air supply pipe through a valve, said lid member being provided with a dust collection opening, and a dust suction duct for connecting an abrasive wheel cover of the grinder to the dust collection opening. The air supply line to the grinder is connected by a T-joint to the compressed air inlet passage.\nHowever, such conventional scrap collection systems have the following disadvantages. A problem with the first system which simply collects falling scraps is that scraps are often prevented from falling because instead they are unintentionally deposited in a mold or the like, and are responsible for making dents in the products molded thereby, thus resulting in an increase in the number of defective products, a shortening in the service life of tools and a lowering of the efficiency of the system.\nA problem with the second system, namely the dust collector, is that the size of the duct collector becomes very large, so that not only is there a need to provide a space sufficient for its installation, but the complicated equipment also leads to an increased cost. Another problem is that as the sharpness of the processing tools deteriorates, scraps are likely to ascend while sticking to the punches because of a lack of suction force of the dust collector, or holes or openings are clogged by scraps."}
{"text": "1. Field of the Invention\nThe present invention relates to a balanced filter having the function of a balun performing conversion between unbalanced and balanced signals and the function of a filter performing band control, and more particularly to a balanced filter effective in reducing a filter size.\n2. Description of the Related Art\nRadio communication equipment comprises various RF (radio frequency) devices, such as an antenna, a filter, an RF switch, a power amplifier, an RF-IC, and a balun. Among these parts, resonance devices, such as an antenna and a filter, handle an unbalanced signal on the basis of the ground potential, while an RF-IC for producing and processing an RF signal handles a balanced signal. A balun functioning as an unbalance-balance transformer is therefore used when those two types of parts are connected to each other.\nThat type of balun is disclosed, for example, in the following Patent Documents:                Patent Document 1: Japanese Unexamined Patent Application Publication No. 2000-134009        Patent Document 2: Japanese Unexamined Patent Application Publication No. 2001-36310        \nThe baluns disclosed in those Patent Documents are of the type that an unbalanced line and a balanced line are coupled through a coupling line. In the structures of those baluns, as shown in FIG. 3 of Patent Document 2, the unbalanced line and the balanced line are formed on one substrate, and the coupling line is formed on another substrate. The coupling line is laid over both the unbalanced line and the balanced line so that the unbalanced line and the balanced line are coupled to each other.\nIn a coupling mode of the balun thus constructed, as shown in FIG. 8 and explained in paragraph 0016 of Patent Document 2, “an unbalanced signal inputted from an unbalanced signal terminal 3 is propagated in the order of a first coupling line 101, a second coupling line 102, and a third coupling line 103”.\nWith the balun structures disclosed in Patent Documents 1 and 2, however, a resulting frequency characteristic is as shown in FIG. 4 of Patent Document 2. Accordingly, the disclosed structures are usable as a balun, but they have a difficulty in ensuring a band characteristic required for the filter.\nOn the other hand, many balanced filters each comprising a balun and a filter combined into an integral unit have recently been devised with the intent to reduce the size of radio communication equipment. That type of balanced filter is disclosed, for example, in the following Patent Document:                Patent Document 3: Japanese Unexamined Patent Application Publication No. 2003-087008        \nThe balanced filter disclosed in Patent Document 3 has a structure in which a filter and a balun each designed using a ¼-wavelength resonator are combined on a dielectric substrate. A dielectric layer constituting the filter and a dielectric layer constituting the balun are formed one above the other in an integral structure.\nAlso, Patent Document 3 discloses a structure in which a DC power supply layer is formed in the balun, for making the balanced filter adaptable for the case where the RF-IC requires a balanced signal superimposed on a DC component. This structure is intended to realize a further reduction of the filter size.\nHowever, the structure in which a balun section and a filter section are separately formed and integrated together has the problem as follows. When the filter function with a high attenuation is demanded, the filter section is required to have a multistage structure. Therefore, satisfactory flexibility in design cannot be ensured in a limited space, and a reduction of the size is very difficult to realize."}
{"text": "Field of the Invention\nThe present invention relates to a scintillator which emits light when radiation is applied thereto, and a radiation detector using the scintillator.\nDescription of the Related Art\nIn a flat panel detector (FPD) used in clinical practice or the like, X-rays which pass through a subject are received by a scintillator, and light emitted from the scintillator is detected by light receiving elements serving as a detector. The light receiving elements are arranged as a two-dimensional array. Japanese Patent Application Laid-Open No. 2001-58881 discloses a scintillator which exhibits an optical waveguide property in order to prevent light emitted by the scintillator from spreading to an adjacent light receiving element under such circumstances. Japanese Patent Application Laid-Open No. 2001-58881 discloses a technology in which, after sacrificial fibers aligned in a direction in which light is sought to he waveguided and scintillator powder are sintered, the sacrificial fibers are removed to obtain a porous scintillator, and pores therein are filled with an absorbing member or the like."}
{"text": "The invention relates to semiconductor memory devices and, more particularly, to a method for fabricating a non-volatile memory device with a charge trapping layer.\nNon-volatile memory devices, which are electrically programmable and erasable devices, are used mainly in electronic components that require data retention even when no power is supplied. Such non-volatile memory devices typically have a floating gate structure. According to the existence or nonexistence of charges in a floating gate, a program operation or an erase operation is performed. As the degree of integration of semiconductor memory devices has increased, non-volatile memory devices with new cell structures have been required. Such non-volatile memory devices with new cell structures include non-volatile memory devices with a charge trapping layer. Such non-volatile memory devices with a charge trapping layer are classified as silicon-oxide-nitride-oxide-silicon (SONOS) structures and metal-aluminum oxide-nitride-oxide-silicon (MANOS) structures, depending on the characteristics of the charge trapping layer. According to a bias voltage applied to non-volatile memory devices with such structures, charges are stored in or discharged from the charge trapping layer to electrically perform a program operation or an erase operation.\nFIG. 1 illustrates a cross-sectional view of a conventional non-volatile memory device with a charge trapping layer 110. Referring to FIG. 1, the non-volatile memory device with the charge trapping layer 110 includes a tunneling layer 105, the charge trapping layer 110, a blocking layer 115, and a control gate electrode 120 stacked on a semiconductor substrate 100. The blocking layer 115 is disposed between the control gate electrode 120 and the charge trapping layer 110 and generally comprises an aluminum oxide (Al2O3) film. Since the quality of the aluminum oxide (Al2O3) film is not high and a grain size of the aluminum oxide (Al2O3) film is small, a leakage current may occur. When the non-volatile memory device is fabricated with the quality of the film not being sufficiently high, characteristics of data retention may be adversely affected. To address this limitation, a thermal treatment process is performed to increase (i.e., “densify”) the quality of the film of the blocking layer 115, after the blocking layer 115 is deposited. A method for the film of the blocking layer 115 includes a process of depositing the blocking layer 115 at a high temperature ranging from approximately 600° C. to approximately 800° C., or a thermal treatment process performed at a high temperature after a deposition process. However, the thermal treatment process performed at the high temperature adversely affects a layer deposited below the blocking layer 115 or a profile of an impurity region formed in the semiconductor substrate 100. Also, temperature and run time of the thermal treatment process are limited to obtain the desired high-quality of the blocking layer 115 using the process of depositing the blocking layer 115 at the high temperature or the thermal treatment process after the deposition process. To address these limitations, a method is required, which can improve the characteristics of the film quality of the blocking layer 115 without affecting a layer below the blocking layer 115 or the semiconductor substrate 100, thereby improving the characteristics of the non-volatile memory device."}
{"text": "In a MOS transistor that amplifies radio frequency signals in, for example, WiMAX, a reduction in the size of each element is required. Thus, for example, the width and length of the gate of an amplifying MOS transistor are subject to variation during manufacturing. Thus, when MOS transistors are cascaded, and amplification is performed, necessary gain may not be achieved under the influence of a variation in gain. Gain “gain” of a MOS transistor can be expressed by equation 1 described below. In equation 1, R represents resistance associated with a MOS transistor, and gm represents mutual conductance.gain=gm·R  [equation 1]\nAccording to the foregoing equation, the gain of a MOS transistor may be maintained constant by compensating for the mutual conductance gm. Thus, hitherto, a mutual conductance compensation circuit that compensates for the mutual conductance of a MOS transistor has been provided in an amplifier circuit.\nFIG. 1 indicates the components of an amplifier circuit provided with mutual conductance compensation circuits. The amplifier circuit includes an amplifier unit 9, a first mutual conductance compensation circuit 7, a second mutual conductance compensation circuit 8, and MOS transistors 4, 5, and 6 for connecting the circuits to each other.\nThe amplifier unit 9 includes a first MOS transistor 1, a second MOS transistor 2, a third MOS transistor 3, a first inductor 15, and a second inductor 16. One end of each of the first inductor 15 and the second inductor 16 is connected to a drain-side voltage source AVD. The other ends of the first inductor 15 and the second inductor 16 are connected to the drains of the second MOS transistor 2 and the third MOS transistor 3, respectively. The source of each of the second MOS transistor 2 and the third MOS transistor 3 is connected to the drain of the first MOS transistor 1. Input signals are input to the gate of each of the second MOS transistor 2 and the third MOS transistor 3. The source of the first MOS transistor 1 is connected to a ground AVS.\nOne end of the first mutual conductance compensation circuit 7 is connected to the drain-side voltage source AVD. The first mutual conductance compensation circuit 7 is a circuit that generates current for controlling mutual conductance so that the mutual conductance is maintained constant. The other end of the first mutual conductance compensation circuit 7 is connected to the drain of the fourth MOS transistor 4. The source of the fourth MOS transistor 4 is connected to the ground AVS. The gate of the fourth MOS transistor 4 is connected to the drain of the MOS transistor 4. Moreover, the gate of the fourth MOS transistor 4 is connected to the gate of the first MOS transistor 1.\nOne end of the second mutual conductance compensation circuit 8 is connected to the drain-side voltage source AVD. The second mutual conductance compensation circuit 8 is a circuit that generates current for controlling mutual conductance so that the mutual conductance is maintained constant. The other end of the second mutual conductance compensation circuit 8 is connected to the drain of the fifth MOS transistor 5. The source of the fifth MOS transistor 5 is connected to the drain of the sixth MOS transistor 6. The source of the sixth MOS transistor 6 is connected to the ground AVS. The gate of the sixth MOS transistor 6 is connected to the gate of the first MOS transistor 1 and the gate of the fourth MOS transistor. The gate of the fifth MOS transistor 5 is connected to the drain of the fifth MOS transistor 5. Moreover, the gate of the fifth MOS transistor 5 is connected to the respective gates of the second MOS transistor 2 and the third MOS transistor 3 via resistors 17 and 18.\nThe first mutual conductance compensation circuit 7 generates current such that the mutual conductance gm of the first MOS transistor 1 is maintained constant and mirrors bias to the first MOS transistor 1, using the fourth MOS transistor 4.\nThe second mutual conductance compensation circuit 8 generates current such that the mutual conductance gm of the second MOS transistor 2 and the third MOS transistor 3 is maintained constant and mirrors bias to the second MOS transistor 2 and the third MOS transistor 3, using the fifth MOS transistor 5.\nAs a MOS transistor becomes smaller, manufacturing errors in the width and length of the gate of a MOS transistor become large, resulting in a difference in the mutual conductance gm, as indicated in FIG. 2. FIG. 2 indicates characteristics of a circuit in FIG. 1. In FIG. 2, a maximum value characteristic Max represents a MOS transistor having the largest mutual conductance. A minimum value characteristic Min represents a MOS transistor having the smallest mutual conductance. Reference letter Typ represents a MOS transistor having mutual conductance of a standard value characteristic. The standard value characteristic Typ is a characteristic based on the gate electrode width and the gate electrode length set in the design stage. It is assumed here that overdrive voltage Vod=gate-source voltage Vgs−threshold voltage Vth.\nIn FIG. 2, in order to compensate for the gain of a MOS transistor the mutual conductance gm of which is lower than that in the standard value characteristic Typ, as indicated by the minimum value characteristic Min, the overdrive voltage Vod needs to be increased. In order to compensate for the mutual conductance gm of each MOS transistor only by bias, Vgs needs to be increased until Vod=Vod0. However, even when Vgs is increased until Vod=Vod0, as indicated in FIG. 2, in a MOS transistor having the minimum value characteristic Min, in which the mutual conductance gm is far less than that in the standard value characteristic Typ, Vgs cannot be increased until the mutual conductance gm reaches mutual conductance gm1 in the standard value characteristic Typ. Thus, when the mutual conductance gm decreases, it may be impossible to compensate for the gain only by controlling the gate-source voltage Vgs of a MOS transistor.\nOn the other hand, in the case of a MOS transistor the mutual conductance gm of which is higher than that in the standard value characteristic Typ, as indicated by the maximum value characteristic Max, in order to compensate for the gain, the gate-source voltage Vgs needs to be controlled to be decreased so that Vod=Vod2. In this case, the mutual conductance gm can be compensated for. However, the overdrive voltage Vod decreases, and thus the amplitude of an input signal that can be linearly amplified decreases. That is, the linearity is deteriorated by compensating for the mutual conductance gm.\nWhen variations occur in elements under the influence of the manufacturing process, the compatibility of compensation for the gain of a MOS transistor with compensation for the linearity cannot be achieved only by controlling the gate-source voltage Vgs.\nJapanese Laid-open Patent Publication No. 2000-174568 is known as a technique for bias control for a variation in the gain of the amplifier unit 9."}
{"text": "Recent advances in both solid-state lighting and rechargeable batteries, has enabled the generation of energy-efficient, compact portable lighting devices that people can carry with them. A portable lighting device is a lighting device that can provide its illumination function without a wired power supply, typically being battery powered by means of a battery on board the lighting device, or potentially instead being powered by another type of on-board power supply such as a manual dynamo, or even being powered using a wireless power transfer (WPT) technique based on radiative electromagnetic induction. A portable lighting device can thus be taken by the user from one location to the next. For instance, a portable lighting device may be taken from a living room to a garden terrace or to a bathroom to create a cozy or romantic ambience.\nTypically the light emitted by known portable lighting devices can be controlled using one or more switches located on the portable lighting. Other portable lighting devices include touch-sensitive surfaces that may be at least partially capacitive. A user's touch may be detected based on a change in that capacitance, and light emitted by one or more light sources of the portable lighting device may be altered based on the nature of the user's touch."}
{"text": "1. Field of the Invention\nThis invention relates to semi-polar optoelectronic devices and methods of fabrication thereof, and in particular, semi-polar III-nitride optoelectronic devices on m-plane substrates with miscuts less than +/−15 degrees in the c-direction.\n2. Description of the Related Art\n(Note: This application references a number of different publications as indicated throughout the specification by one or more reference numbers within brackets, e.g., [x]. A list of these different publications ordered according to these reference numbers can be found below in the section entitled “References.” Each of these publications is incorporated by reference herein.)\nThe anticipated high commercial demand for next-generation display technologies, such as miniature mobile projectors and high-definition fly-spot displays, has significantly accelerated the development of direct-emission green laser diodes (LDs). Technical criteria for such applications require LDs to have high efficiency, reliability, compactness, and modulation response capabilities [1]. While the wurtzite (Al,Ga,In)N-based material system is largely agreed upon as the leading candidate for green optoelectronic devices, a general consensus behind which crystal plane is optimal for epitaxial growth has yet to be reached.\nContinuous-wave (CW) operation of LDs in the green spectral region has been demonstrated for devices grown on the conventional c-plane of GaN [2-4]. These devices, however, suffer from parasitic internal electric fields that give rise to the quantum-confined stark effect (QCSE) which decreases the quantum-well (QW) radiative recombination rate, and causes a blue-shift in emission wavelength with increasing carrier injection [5]. In addition, QCSE becomes more significant for long wavelength optoelectronic devices due to the increased lattice mismatch between the In-rich InGaN QWs and the barrier [2].\nTo circumvent polarization effects, researchers have demonstrated operation of LDs grown on the nonpolar m-plane orientation of the wurtzite crystal [6]. While a promising candidate for high power blue LDs, m-plane LDs has so far been limited to 500 nm lasing emission due to the formation of stacking faults in the active region [7-11].\nVarious semi-polar (or semipolar) planes, such as (10-11) and (11-22), have been also investigated as alternative growth planes [12-13]. Recently, researchers have reported lasing in the green spectral region from high quality InGaN QWs grown on the semi-polar (20-21) plane [14-15]. Further studies have shown that green emitting QWs grown on (20-21) exhibit high compositional homogeneity with localization energy values lower than that reported for c-plane [16]. However, quaternary AlInGaN cladding layers were needed in (20-21) LDs to achieve sufficient modal confinement without generating strain-induced misfit dislocations (MDs) [17]. While the use of quaternary AlInGaN cladding layers presents a solution for semi-polar planes with low critical thickness, a simple AlGaN cladding-free structure with high composition InGaN waveguides would be much more attractive from a mass-production point of view [7, 15, 18].\nWhat is needed are improved methods for optoelectronic device growth. The present invention satisfies this need."}
{"text": "Users of a multimedia content distribution network (MCDN) may be provided a wide range of video assets to select from. A service provider operating the MCDN may be faced with various quality control issues related to the video assets and the performance of MCDN equipment. Feedback about MCDN performance issues are typically obtained via information gleaned from user support requests and/or costly support visits to user locations."}
{"text": "The present invention relates to a method for manufacturing a dynamic quantity detection device that is formed by bonding a semiconductor chip that includes a detection element for detecting a dynamic quantity to a stand using a bonding layer.\nThe publication JP-A-2000-241273 discloses such a method in the manufacturing process of a discrete pressure detection device 100 shown in FIGS. 3A and 3B.\nAs shown in FIGS. 3A and 3B, the discrete pressure detection device 100 includes a metal stem 20, a discrete semiconductor chip 30, and a glass layer 40. The metal stem 20 includes a disk-like diaphragm 10. The semiconductor chip 30 is bonded to the diaphragm 10 with the glass layer 40. The discrete semiconductor chip 30 is in contact with the glass layer 40 at a first surface 31 of the discrete semiconductor chip 30. The diaphragm 10 is in contact with the glass layer 40 at a first surface 11 of the diaphragm 10. The discrete semiconductor chip 30 includes four gauges (detection elements) 51, 52, 53, 54, or four piezo resistors 51, 52, 53, 54. The gauges 51, 52, 53, 54 are located in a second surface 32 of the chip 30, which is opposite to the first surface 31. The glass layer 40 is, for example, made of a low melting point glass.\nWhen a medium that transmits pressure to be detected is introduced into the stem 20 to apply the pressure to a second surface 12 of the diaphragm 10, which is opposite to the first surface 11, the diaphragm 10 strains and the discrete semiconductor chip 30 also strains synchronously to change the resistances of the gauges 51, 52, 53, 54. The resistance changes are processed by a processing circuit, which is not shown in the figure, in order to detect the pressure. The processing circuit is provided in another chip outside the discrete semiconductor chip 30 and includes processing circuit elements such as a bipolar transistor and a trimming resistor.\nBecause the processing circuit is located in another chip, the system that includes the discrete pressure detection device 100 and the chip including the processing circuit is relatively bulky. To address the issue of the bulkiness of the system, the inventors of the present invention attempted to integrate gauges 51, 52, 53, 54 and processing circuit elements for making up a processing circuit equivalent to the foregoing processing circuit in an intelligent semiconductor chip 33 and to form an intelligent pressure detection device 101, which is shown in FIGS. 4A and 4B. The processing circuit elements are included in a processing circuit area 70.\nHowever, it turned out that the following issue was caused when the intelligent pressure detection device 101 was manufactured using the intelligent semiconductor chip 33. When the discrete pressure detection device 100 of FIGS. 3A and 3B is manufactured, the discrete semiconductor chip 30 is normally bonded to the metal stem 20 by sintering the glass layer 40 at a temperature of approximately 500° C. However, when the glass layer 40 was sintered at a temperature of approximately 500° C. the characteristics of a bipolar transistor, which is included in the processing circuit area 70 as a processing circuit element, changed.\nThe characteristics change of the bipolar transistor may be caused by the following mechanism. Although not shown in the figure, the intelligent semiconductor chip 33 includes a substrate, which is made of silicon. Wiring lines, which electrically interconnect elements such as a bipolar transistor located in a surface of the substrate, an oxide film, which insulates the wiring lines, and a passivation film, which protects the wiring lines and the oxide film, are located on the surface of the substrate. The oxide film and the passivation film are formed in the atmosphere that contains hydrogen in the wafer manufacturing process, in which a wafer is processed for manufacturing the intelligent semiconductor chip 33.\nWhen the oxide film and the passivation film are formed, the hydrogen in the atmosphere is absorbed into the oxide film and the passivation film as hydrogen ions. The dangling bonds of the silicon atoms located at the interface between the surface of the substrate and the oxide film are terminated by the hydrogen ions. At a subsequent step in the wafer manufacturing process, the wafer is heated at a high temperature in the atmosphere that contains hydrogen. During the step, the amount of hydrogen ions in the oxide film and the passivation film can increase as hydrogen in the atmosphere is absorbed, and the bonds between the hydrogen-terminated silicon atoms and the terminating hydrogen atoms at the interface between the surface of the substrate and the oxide film may be broken due to the high temperature.\nAt a later step, the wafer is diced to form the intelligent semiconductor chip 33. Because no heat is applied to the intelligent semiconductor chip 33 during the dicing step, the characteristics of the bi-polar transistor such as current amplification factor in the wafer manufacturing process is substantially determined by the concentration of the hydrogen ions resulting from the above high temperature treatment.\nHowever, when the glass layer 40 is sintered for bonding the intelligent semiconductor chip 33 to the metal stem 20, the processing circuit area 70 is heated. As a result, the hydrogen atoms bonded to the silicon atoms at the interface between the surface of the substrate and the oxide film move away from the silicon atoms, and the hydrogen ions absorbed in the oxide film migrate toward the passivation film or get outgassed into the atmosphere.\nAs a result, electrons trapped by the hydrogen atoms included in the oxide film are released above the surface of the bipolar transistor located in the processing circuit area 70, and the transistor injection efficiency increases. Consequently, the base current increases, and the collector current also increases among the currents that flow through the base, the emitter, and the collector of the bipolar transistor. Thus, the current amplification factor, which is one of the characteristics of the bipolar transistor, increases."}
{"text": "This invention relates to an unloader for a tower grain dryer or the like, and more particularly to such an unloader which meters or positively unloads the dried grain from such a dryer in such manner that dried grain is substantially uniformly removed from all regions of the dryer outlet section.\nTower grain dryers are well known. Generally, they comprise a vertical tower of cylindrical or other shape. The tower has a plenum located within the tower and spaced from the inner surfaces of the outer walls of the tower. Typically, the walls of the tower and of the plenum are of porous construction such that air from within the plenum may be forced through the walls. Grain to be dried is conveyed to the upper reaches of the tower and fills the space between the outer tower wall and the plenum so as to form a drying path for the grain to be dried. As dried grain is continuously discharged from the lower end of the grain drying path, additional grain to be dried is loaded into the upper end of the drying path.\nThe plenum is supplied with heated air by means of one or more burner/blower assemblies. The heated air is forced from within the plenum through the porous or perforated walls of the plenum, through the grain in the drying path, and is forced through the porous outer wall of the tower carrying away moisture from the grain.\nReference may be made to the following U.S. Pat. Nos. which describe tower or other similar types of grain dryers of the type discussed above: 746,410, 3,333,348, 3,710,449, 3,766,664, 3,864,845, 3,896,562, 3,955,288, 4,004,351, 4,398,356, 4,423,557, 4,914,834, 5,129,164, and 5,136,791.\nIn certain of these prior art tower grain dryers, the lower portion of the drying path converges inwardly of the tower in a generally conical shape and is divided into a number of separate converging channels. The lower ends of the channels discharge the dried grain into the bottom of the tower to be unloaded. Unloading can be accomplished in a variety of ways. Of course, grain being a fluent material may flow by gravity out through gravity outlets. Alternatively, the grain may be swept from the base of the tower by a sweep auger or the like and conveyed by the sweep auger to a discharge outlet, as described in the above-noted U.S. Pat. No. 3,896,562. Still further, discharge augers may be located directly below the grain drying paths, as shown in U.S. Pat. No. 3,864,845, to unload the dried grain. In U.S. Pat. No. 3,333,348, a grain drying apparatus is disclosed having an inner wall and a spaced outer wall defining a grain flow path and having a pair of rotary agitator arms which extend out from the center of the dryer. In U.S. Pat. No. 4,004,351 uses rotary driven scoop arms to positively remove grain from the grain flow path.\nHowever, it has been found that with tower dryers in which the grain drying path converges inwardly at the bottom of the tower and in which the grain drying path is divided into a number of converging channels, dried grain is oftentimes not uniformly removed from all of the channels. It has been found that grain will be removed at faster rates from some of the channels than from others. Because the speed at which the grain moves through the drying path (and hence the amount of time the grain is exposed to the drying air) varies in direct relation to the rate at which the dried grain is removed from the outlets of the channels, prior grain unloaders which did not uniformly remove dried grain from all of the grain channels caused some of the grain to move too fast through the dryer such that this faster moving grain might not be sufficiently dried, while other grain in the drying path might move too slowly such that this slower moving grain might become over dried. Of course, under dried grain is not desirable because it could lead to spoilage of the grain, and over dried grain is not desirable because it might become damaged from excessive exposure to the heated air which would waste energy and may damage the grain.\nIn the grain dryer described in the parent application of the present invention, now U.S. Pat. No. 5,566,470, dried grain is removed from the grain flow path by a rotary driven metering unloader body rotatably driven about a vertical axis and having a vertical wall for blocking the flow of grain from the grain flow path with the unloader wall having a grain inlet opening such that grain may enter the unloader body only by way of the grain inlet opening. Upon rotation of the unloader body, the grain within the unloader body is conveyed toward a center grain discharge outlet opening for removal from the dryer. While this unloader works well for its intended purpose, in order for it to vary the rate at which dried grain is unloaded, the rotary speed at which the unloader body is driven must be varied. This requires the use of a variable speed drive which is costly. In addition, the earlier metering unloader of this invention required a sweeping action to move the grain across the horizontal floor of the unloader to the grain discharge outlet. Therefore, there has been a need for a metering unloader body which does not require the use of variable speed drive and which does not require a sweeping action to discharge the grain."}
{"text": "This invention is concerned with the separation of a mixture of carbon dioxide and ethane derived from a prior separation or recovery stage or process. More particularly, the invention pertains to separation of a carbon dioxide-ethane azeotrope by freezing the carbon dioxide and removing the ethane with propane or a heavier hydrocarbon.\nIt is sometimes desirable to separate carbon dioxide from a gaseous mixture containing ethane, for example, in the treatment of natural gas or a reservoir injection gas, or in the preparation of synthetic natural gas. In one or more stages of these processes, carbon dioxide, ethane and heavier hydrocarbons are separated or recovered as a mixture. It is then desirable to separate the carbon dioxide from the remaining hydrocarbons, sometimes called NGL (natural gas liquids). When carbon dioxide is distilled or fractionated from the hydrocarbons, an azeotrope of carbon dioxide and ethane is formed at an overhead point in the distillation column. At this point no more ethane is separated from the carbon dioxide. Generally, unless it is desirable to recover the carbon dioxide, the carbon dioxide is vented unless the concentration of ethane in the azeotrope is such that environmental regulations make it necessary to burn the ethane-carbon dioxide mixture. Flaring is frequently an expensive and difficult process, especially in areas where the amount and concentration of the gas being flared randomly varies with various operating or producing conditions. In processes where it is desirable to recover the carbon dioxide, for example, tertiary recovery processes, the unseparated ethane is lost.\nThe ethane in the carbon dioxide is a valuable hydrocarbon. This invention is concerned with separating and recovering ethane along with the other hydrocarbons that might be in a carbon dioxide-ethane mixture derived from a prior process."}
{"text": "In the manufacture of transistors in a semiconductor substrate, such as a silicon wafer, several steps are required to form the gates of the transistors. A first layer of a gate oxide, e.g., silicon oxide, is deposited over the silicon substrate to protect it. Next a conductive, doped polysilicon layer is deposited, optionally followed by a layer of a conductive material such as a refractory metal silicide. A cap layer of silicon oxide or silicon nitride, preferably a silicon oxide produced by the chemical vapor deposition from tetraethoxysilane (TEOS), is also deposited.\nA typical gate stack is shown in FIG. 1A wherein a silicon wafer 12 has a layer of gate oxide 13 thereon, followed by a first gate layer of conductive, highly doped polysilicon 14. A second gate layer 15 of a refractory metal silicide, such as titanium silicide or tungsten silicide, and a cap layer 16 of silicon oxide are also deposited.\nThe gate stack must now be etched to form a gate between the areas in the substrate in which the source and drain will be formed. The silicon oxide hard mask layer 16 is coated with photoresist, the photoresist is patterned, and the silicon oxide layer etched in conventional manner to form a \"hard mask\". Suitably the etchant for the hard mask is a fluorine-containing gas such as CHF.sub.3, C.sub.2 F.sub.6, CF.sub.4 and the like.\nFIG. 1B illustrates the gate stack after it has been etched, the same numerals used as in FIG. 1A for the same layers. The silicon oxide hard mask and the gate oxide layer remain, and the two conductive layers are etched.\nAs devices become smaller and are placed closer together, the gates become correspondingly thinner and narrower. Thus wet etch techniques have been largely replaced by dry etching, particularly plasma etching, because plasma etching is more anisotropic and thus does not undercut the gate, particularly at its intersection with the substrate. Also, plasma etchants can be used that are highly selective; that is, they etch certain materials faster than others. In the case of a gate stack, the etchant must etch the silicide layer and the polysilicon layer, but not silicon oxide; thus the cap layers and the initial gate oxide layers remain intact, as shown in FIG. 1B. Plasma processing is also favored because processing can be carried out at comparatively low temperatures that will not damage the substrate or previously formed devices.\nVarious vacuum chambers have been designed to carry out plasma etching. One particular design has been disclosed by Ogle in U.S. Pat. No. 4,948,458. The etch apparatus forms a magnetically coupled planar plasma for treatment of a substrate such as a silicon wafer. This apparatus comprises a chamber having a dielectric window. A planar coil is mounted proximate to the window, and an RF power source is coupled to the coil, generally through an impedance matching circuit to maximize power transfer, and a timing circuit to provide resonance at the operating frequency, typically 13.56 MHz. Process gas is supplied to the chamber through an inlet port. When the RF power is turned on, a planar magnetic field is induced which extends into the interior of the chamber through the dielectric window. Thus a circulating flow of electrons is induced in the chamber in a plane parallel to the planar coil, limiting the transfer of kinetic energy in non-planar directions.\nA surface support in the chamber supports the substrate to be processed parallel to the plane of the coil, and thus to the plasma. Since the plasma species have little velocity in non-planar directions, the kinetic impact of ionic species on the substrate is small and the predominant reaction is a chemical reaction between the plasma species and the layers on the substrate. This reaction takes place at low temperatures and low pressures.\nA detailed description of the above apparatus is shown in FIGS. 2-4. Referring to FIGS. 2, 3 and 4, a plasma treatment system 100 suitable for etching individual semiconductor wafers 110, includes a chamber 112 having an access port 114 formed in an upper wall 116. A dielectric window 118 is disposed below the upper wall 116 and extends across the access port 114. The dielectric window 118 is sealed to the wall 116 to define a vacuum sealed interior 119 of the chamber 112. A planar coil 120 is mounted adjacent to the dielectric window 118. The coil 120 is formed as a spiral having a center tap 122 and an outer tap 124. The plane of the coil 120 is oriented parallel both to the dielectric window 118 and to the support surface 113 upon which the wafer 110 is mounted during etching. The coil 120 is able to produce a planar plasma within the interior 119 of the chamber 112 which is parallel to the wafer 110. A suitable distance between the coil 120 and the support surface 113 can be about 5-10 cm.\nA radiofrequency (RF) generator 130 leads to a coaxial cable 132 to a matching circuit 134. The matching circuit 134 includes a primary coil 136 and a secondary loop 138 that may be positioned to adjust the effective coupling of the circuit and allow loading of the circuit at the frequency of operation. The primary coil 136 may be mounted on a disk 140 that can be rotated about a vertical axis 142 to adjust the coupling. A variable capacitor 144 is in series with the secondary loop 138 to adjust the circuit resonant frequency with the frequency output of the RF generator 130. Impedance matching maximizes the efficiency of the power transfer to the planar coil 120. An additional capacitor 146 is provided in the primary circuit to cancel part of the inductive reactance of the primary coil 136 in the circuit.\nA second RF potential from a source 172 is applied to the wafer support 113. This RF generator 172 can be operated at a low frequency (below about 550 kHz) or a high frequency (13.56 MHz) but generally a low frequency is employed. The frequency of the RF generator 172 and that of the first RF generator 130 is generally different, and provides a resonant current flow in the coil 120. Thus if the primary RF generator 130 operates at 13.56 MHz, the second RF generator 172 operates suitably at 400 kHz for example. This two RF power source configuration allows control of the amount of energy introduced to the system 100 through the primary RF generator 130 and control of the power output of the RF generator 172 permits control of the velocity imparted to the reactive species in the plasma.\nIn order to maximize the density of the plasma, and therefore to reduce the time required to perform an etch step, comparatively high power is applied to the planar coil 120, on the order of 300-600 watts. A lower power is applied to the substrate support 113 of about 75-300 watts.\nThus high power provides fast etch rates at a low pressure in the chamber, i.e., 1 Torr or less. Further, since the kinetic energy of ions in such a plasma is low, there is less damage to the wafer.\nHowever, we have found that charging of the wafers processed in the above planar coil apparatus is a problem. A net charge on the substrate support 113 is transferred to the wafer 110 being processed. It is believed this net charge is caused by non-uniformities in the plasma. Further, the selectivity of halogen-containing etchants, such as chlorine, is reduced by a high ion density in the plasma. The high density plasma can also result in undercutting of the wafer stack to be etched. For example, a high density plasma contains, of necessity, a high density of radical species that etch laterally as well as horizontally, causing undercutting of the gates, as shown in FIG. 5, discussed further hereinbelow. In order to reduce this undercutting, it has become common to provide sidewall protection of the gate during further plasma processing. During the course of the etching, to prevent lateral etching and undercutting of the etched gate stack, particularly at the base of the gate stack, it is also known to generate a plasma that, in addition to etching the gate stack, also deposits a protective film on the freshly etched sidewall surfaces. The formation of such layers, which can be polymeric, for example, improves the anisotropy of the etch and prevents undercutting. However, the thickness of this sidewall deposit is difficult to control; if it is too thin, it does not protect the etched sidewalls; and if it is too thick, it takes up more space on the substrate, and impairs the critical dimensions of the gates and their devices. This is unacceptable for submicron design rules.\nTo perform an etch step in the apparatus of FIGS. 2-4, a process gas is introduced into the chamber 112 through a port 150 formed in the side of the chamber 110. A vacuum exhaust system 152 maintains the pressure in the system and evacuates volatile by-product and unreacted gases.\nThus an improved method of etching gate stacks having good etch rates and good selectivity but improved anisotropy, would be highly desirable."}
{"text": "Balloon angioplasty is a widely used procedure for expanding constricted body passageways, such as arteries and other blood vessels, or various ducts (e.g., of the biliary system). In an angioplasty procedure, an uninflated angioplasty balloon attached to a catheter is delivered to a constricted region of a body passageway. Once the balloon is in position at the constricted region, fluid is injected through a lumen of the catheter and into the balloon. The balloon consequently inflates and exerts pressure against the constricted region to expand the passageway. After use, the balloon is collapsed, and the catheter is withdrawn. Although treatment of constricted arteries in the vasculature is one common example where balloon catheters are used, this is only one example of how balloon catheters may be used and many other uses are possible. For example, balloon catheters may also be used to temporarily occlude vessels during medical procedures to prevent blood or other fluids from flowing through a vessel. Balloon catheters may also be used to expand other intraluminal devices without dilating the surrounding vessel wall, such as stent-grafts that may be used to treat aneurysms. The above-described examples are only some of the applications in which balloon catheters are used by physicians. Many other applications for balloon catheters are known and/or may be developed in the future.\nBalloons have a number of important design parameters. One is rated burst pressure, which is the statistically-determined maximum pressure to which a balloon may be inflated without rupturing. In order to expand hard, calcified lesions, it is desirable that the balloon have a rated burst pressure of at least 15 atm. It is also desirable that the balloon have a low wall thickness to minimize the profile of the delivery system. A wall thickness of about 0.03 millimeters or lower is generally preferred. For a given balloon material, however, there is a trade-off between burst pressure and wall thickness, in that the burst pressure generally decreases when the wall thickness is reduced. Accordingly, there is a need for a means of increasing the strength of balloon materials to attain higher rated burst pressures at lower wall thicknesses."}
{"text": "DNA microarrays have been used to measure gene expression in samples from patients and to facilitate diagnosis of disease. Gene expression can reveal the presence of cancer in a patient in addition to the type, stage, and origin of the cancer. Gene expression may even have a role in predicting the efficacy of cancer therapies. In recent decades, the National Cancer Institute (NCI) has tested cancer therapeutics for their effect in limiting the growth of 60 human cancer cell lines. The NCI has also measured gene expression in those 60 cancer cell lines using DNA microarrays. Various studies have explored the relationship between gene expression and therapeutic effect using the NCI datasets.\nDuring cancer treatment, critical time is often lost due to a trial and error approach to finding an effective therapy. In addition, cancer cells often develop resistance to a previously effective therapy. In such situations, patient outcome would be greatly improved by early detection of such resistance.\nThus, there exists a need in the art for methods and devices that can predict the responsiveness of cancer patients to a medical treatment."}
{"text": "1. Field of Invention\nThis invention relates to molecules useful in determining minimal functional regions of IGFBP-3 and for antagonizing an IGF-I or IGF-II activity.\n2. Description of Background and Related Art\nThe insulin-like growth factors I and II (IGF-I and IGF-II, respectively) mediate multiple effects in vivo, including cell proliferation, cell differentiation, inhibition of cell death, and insulin-like activity (reviewed in Clark and Robinson, Cytokine Growth Factor Rev., 7: 65–80 (1996); Jones and Clemmons, Endocr. Rev., 16: 3–34 (1995)). Most of these mitogenic and metabolic responses are initiated by activation of the IGF-I receptor, an α2β2-heterotetramer closely related to the insulin receptor (McInnes and Sykes, Biopoly., 43: 339–366 (1998); Ullrich et al., EMBO J., 5: 2503–2512 (1986)). The IGF-I and insulin receptors bind their specific ligands with nanomolar affinity. IGF-I and insulin can cross-react with their respective non-cognate receptors, albeit at a 100–1000-fold lower affinity (Jones and Clemmons, supra). The crystal structure describing part of the extracellular portion of the IGF-I receptor has been reported (Garrett et al., Nature, 394: 395–399 (1998)).\nUnlike insulin, the activity and half-life of IGF-I are modulated by six IGF-I binding proteins (IGFBPs 1–6), and perhaps additionally by a more distantly related class of proteins (Jones and Clemmons, supra; Baxter et al., Endocrinology, 139: 4036 (1998)). IGFBPs can either inhibit or potentiate IGF activity, depending on whether they are soluble or cell-membrane associated (Bach and Rechler, Diabetes Reviews,3: 38–61 (1995)). The IGFBPs bind IGF-I and IGF-II with varying affinities and specificities (Jones and Clemmons, supra; Bach and Rechler, supra). For example, IGFBP-3 binds IGF-I and IGF-II with a similar affinity, whereas IGFBP-2 and IGFBP-6 bind IGF-II with a much higher affinity than they bind IGF-I (Bach and Rechler, supra; Oh et al., Endocrinology, 132: 1337–1344 (1993)). WO 01/75064 discloses additional human secreted IGFBP-like polypeptides that are encoded by nucleic acid sequences isolated from cDNA libraries from adrenal gland mRNA and thymus mRNA.\nStructurally, IGF-I is a single-chain, 70-amino-acid protein with high homology to proinsulin. Unlike the other members of the insulin superfamily, the C region of the IGF's is not proteolytically removed after translation. The solution NMR structures of IGF-I (Cooke et al., Biochemistry, 30: 5484–5491 (1991); Hua et al., J. Mol. Biol., 259: 297–313 (1996)), mini-IGF-I (an engineered variant lacking the C-chain; DeWolf et al., Protein Science, 5: 2193–2202 (1996)),and IGF-II (Terasawa et al., EMBO J., 13: 5590–5597 (1994); Torres et al., J. Mol. Biol., 248: 385–401 (1995)) have been reported. It is generally accepted that distinct epitopes on IGF-I are used to bind receptor and binding proteins. It has been demonstrated in animal models that receptor-inactive IGF mutants are able to displace endogenous IGF-I from binding proteins and thereby generate a net IGF-I effect in vivo (Loddick et al., Proc. Natl. Acad. Sci. USA, 95: 1894–1898 (1998); Lowman et al., Biochemistry, 37: 8870–8878 (1998); U.S. Pat. Nos. 6,121,416 and 6,251,865). While residues Y24, Y29, Y31, and Y60 are implicated in receptor binding, IGF mutants thereof still bind to IGFBPs (Bayne et al., J. Biol. Chem., 265: 15648–15652 (1990); Bayne et al., J. Biol. Chem., 264: 11004–11008 (1989); Cascieri et al., Biochemistry, 27: 3229–3233 (1988); Lowman et al., supra).\nAdditionally, a variant designated (1–27,gly4, 38–70)-hIGF-I, wherein residues 28–37 of the C region of human IGF-I are replaced by a four-residue glycine bridge, has been discovered that binds to IGFBP's but not to IGF receptors (Bar et al., Endocrinology, 127: 3243–3245 (1990)).\nA multitude of mutagenesis studies have addressed the characterization of the IGFBP-binding epitope on IGF-I (Bagley et al., Biochem. J., 259: 665–671 (1989); Baxter et al., J. Biol. Chem., 267: 60–65 (1992); Bayne et al., J. Biol. Chem., 263: 6233–6239 (1988); Clemmons et al., J. Biol. Chem., 265: 12210–12216 (1990); Clemmons et al., Endocrinology, 131: 890–895 (1992); Oh et al, supra). In summary, the N-terminal residues 3 and 4 and the helical region comprising residues 8–17 were found to be important for binding to the IGFBPs. Additionally, an epitope involving residues 49–51 in binding to IGFBP-1, -2 and -5 has been identified (Clemmons et al., Endocrinology, supra, 1992). Furthermore, a naturally occurring truncated form of IGF-I lacking the first three N-terminal amino acids (called des (1-3)-IGF-I) was demonstrated to bind IGFBP-3 with 25 times lower affinity (Heding et al., J. Biol. Chem. 271: 13948–13952 (1996); U.S. Pat. Nos. 5,077,276; 5,164,370; and 5,470,828).\nIn an attempt to characterize the binding contributions of exposed amino acid residues in the N-terminal helix, several alanine mutants of IGF-I were constructed (Jansson et al., Biochemistry, 36: 4108–4117 (1997)). However, the circular dichroism spectra of these mutant proteins showed structural changes compared to wild-type IGF-I, making it difficult to clearly assign IGFBP-binding contributions to the mutated side chains. A different approach was taken in a very recent study where the IGFBP-1 binding epitope on IGF-I was probed by heteronuclear NMR spectroscopy (Jansson et al., J. Biol. Chem., 273: 24701–24707 (1998)). The authors additionally identified residues R36, R37 and R50 to be functionally involved in binding to IGFBP-1.\nOther IGF-I variants have been disclosed. For example, in the patent literature, WO 96/33216 describes a truncated variant having residues 1–69 of authentic IGF-I. EP 742,228 discloses two-chain IGF-I superagonists that are derivatives of the naturally occurring single-chain IGF-I having an abbreviated C domain. The IGF-I analogs are of the formula: BCn, A wherein B is the B domain of IGF-I or a functional analog thereof, C is the C domain of IGF-I or a functional analog thereof, n is the number of amino acids in the C domain and is from about 6 to about 12, and A is the A domain of IGF-I or a functional analog thereof.\nAdditionally, Cascieri et al., Biochemistry, 27: 3229–3233 (1988) discloses four mutants of IGF-I, three of which have reduced affinity to the Type 1 IGF receptor. These mutants are: (Phe23, Phe24, Tyr25) IGF-I (which is equipotent to human IGF-I in its affinity to the Types 1 and 2 IGF and insulin receptors), (Leu24) IGF-I and (Ser24) IGF-I (which have a lower affinity than IGF-I to the human placental Type 1 IGF receptor, the placental insulin receptor, and the Type 1 IGF receptor of rat and mouse cells), and desoctapeptide (Leu24) IGF-I (in which the loss of aromaticity at position 24 is combined with the deletion of the carboxyl-terminal D region of hIGF-I, which has lower affinity than (Leu24)IGF-I for the Type 1 receptor and higher affinity for the insulin receptor). These four mutants have normal affinities for human serum binding proteins.\nBayne et al., J. Biol. Chem., 264: 11004–11008 (1988) discloses three structural analogs of IGF-I: (1–62) IGF-I, which lacks the carboxyl-terminal 8-amino-acid D region of IGF-I; (1–27,Gly4, 38–70) IGF-I, in which residues 28–37 of the C region of IGF-I are replaced by a four-residue glycine bridge; and (1–27,Gly4, 38–62) IGF-I, with a C region glycine replacement and a D region deletion. Peterkofsky et al., Endocrinology, 128: 1769–1779 (1991) discloses data using the Gly4 mutant of Bayne et al., supra, Vol. 264. U.S. Pat. No. 5,714,460 refers to using IGF-I or a compound that increases the active concentration of IGF-I to treat neural damage.\nCascieri et al., J. Biol. Chem., 264: 2199–2202 (1989) discloses three IGF-I analogs in which specific residues in the A region of IGF-I are replaced with the corresponding residues in the A chain of insulin. The analogs are: (Ile41,Glu45,Gln46,Thr49,Ser50,Ile51,Ser53,Tyr55,Gln56)IGF-I, an A chain mutant in which residue 41 is changed from threonine to isoleucine and residues 42–56 of the A region are replaced; (Thr49,Ser50,Ile51)IGF-I; and (Tyr55,Gln56)IGF-I.\nSliecker et al., Adv. Experimental Med. Biol., 343: 25–32 (1994)) describes the binding affinity of various IGF and insulin variants to IGFBPs, IGF receptor, and insulin receptor.\nIGFBPs are secreted by cells in culture and either inhibit or enhance IGF-stimulated functions (Clemmons et al., (1991) In Modern Concepts of Insulin-like Growth Factors. E. M. Spencer, editor. Elsevier, New York, N.Y. 475–486). Known forms of IGFBPs include IGFBP-1, having a molecular weight of approximately 30–40 kDa in humans. See, e.g., WO89/09792, published Oct. 19, 1990, pertaining to cDNA sequences and cloning vectors for IGFBP-1 and IGFBP-2; WO89/08667, published Sep. 21, 1989, relating to an amino acid sequence of IGFBP-1; and WO89/09268, published Oct. 5, 1989, relating to a cDNA sequence of IGFBP-1 and methods of expression for IGFBP-1.\nIGFBP-2 has a molecular weight of approximately 33–36 kDa. See, e.g., Binkert et al., The EMBO Journal, 8: 2497–2502 (1989), relating to a nucleotide and deduced amino acid sequence for IGFBP-2.\nIGFBP-3 has a non-glycosylated molecular weight of about 28 kDa. See, e.g., Baxter et al., Biochim. Biophys. Res. Com., 139:1256–1261 (1986), pertaining to a glycosylated 53-kDa subunit of IGFBP-3 that was purified from human serum; Wood et al., Mol. Endocrinol., 2:1176–1185 (1988), relating to a full-length amino acid sequence for IGFBP-3 and cellular expression of the cloned IGFBP-3 CDNA in mammalian tissue culture cells; WO 90/00569, published Jan. 25, 1990, relating to isolating from human plasma an acid-labile subunit (ALS) of IGFBP complex and the particular amino acid sequence for ALS pertaining to a subunit of IGFBP-3; and Schmid et al., Biochim. Biophys. Res Com., 179: 579–585 (1991), relating to effects of full-length and truncated IGFBP-3 on two different osteoblastic cell lines.\nAlthough initially some inconsistencies in nomenclature for IGFBP4, IGFBP-5, and IGFBP-6 existed, in 1991 participants of the 2nd International IGF Symposium agreed upon an accepted IGFBP4, IGFBP-5, and IGFBP-6 nomenclature. Using accepted terminology, Mohan et al., Proc. Natl. Acad. Sci., 86:8338–8342 (1989) relates to an N-terminal amino acid sequence for an IGFBP-4 isolated from medium conditioned by human osteosarcoma cells, and Shimasaki et al., Mol. Endocrinology, 4:1451–1458 (1990) pertains to IGFBP cDNAs encoding IGFBP-4 from rat and human. WO92/03471 published Mar. 5, 1992, relates to an IGFBP4 (originally designated therein as IGFBP-5); and WO92/03470 published Mar. 5, 1992 relates to genetic material encoding IGFBP4 (originally designated therein as IGFBP-5).\nWO92/12243 published Jul. 23, 1992, relates to IGFBP-5 (originally designated therein as IGFBP-6). Andress and Birnbaum, Biochim. Biophys. Res Com., 176: 213–218 (1991) relates to the modulation of cellular action of a mixture of affinity-purified IGFBPs from U-2-cell-conditioned media on IGFs. WO92/03469 published Mar. 5, 1992, relates to genetic material encoding IGFBP-6 (originally designated therein as IGFBP4); and WO92/03152 published Mar. 5, 1992, relates to an IGFBP-6 (originally designated therein as IGFBP-4). See also U.S. Appln. No. 2003/0082744 as well as U.S. Pat Nos. 6,025,465 and 5,212,074, which disclose IGFBP-6 and its fragments.\nZapf et al., J. Biol. Chem., 265:14892–14898 (1990) pertains to four IGFBPs (IGFBP-2, IGFBP-3, a truncated form of IGFBP-3, and IGFBP-4) isolated from adult human serum by insulin-like growth factor (IGF) affinity chromatography and high-performance liquid chromatography. Shimasaki et al., 2nd International IGF Symposium Abstract (January 1991) discusses amino-terminal amino acids for IGFBP4, IGFBP-5, and IGFBP-6.\nWhen administered alone, i.e., without any IGF, the IGFBPs may also be therapeutically useful for blocking the adverse effects of IGFs, such as those that occur when IGFs are produced in excess, e.g. free IGFs secreted by certain cancer cells such as hormone-producing cancer cells such as breast or kidney cancer cells. More recently, it was demonstrated that U-2 human osteosarcoma cells secrete IGFBP-5 and IGFBP-6 (Andress and Birnbaum, supra; Shimasaki et al., J. Biol. Chem. 266: 10646–10653 (1991); Shimasaki et al., Mol. Endocrinol., 5: 938–948 (1991)). Although affinity-purified IGF-binding proteins derived from U-2-conditioned medium clearly enhanced IGF-I stimulated mitogenesis (Andress and Birnbaum, supra), it was unclear from those studies which protein was responsible for this effect. Mohan et al. demonstrated that IGFBP-4, purified from TE-89 human osteosarcoma cells, inhibits IGF-stimulated osteoblast mitogenesis (Mohan et al., Proc. Natl. Acad. Sci. (U.S.A.), 86: 8338–8342 (1989); see also LaTour et al., Mol. Endocrinol., 4: 1806–1814 (1990)). WO 03/068160 discloses use of IGFBP-3 for inhibiting tumor growth. WO 03/006029 discloses a method of inducing apoptosis in a cancer cell comprising increasing the expression of IGFBP-5 by the cell to an apoptosis-inducing amount. A method of killing cancer cells, a method of sensitizing cancer cells to agents that induce apoptosis, and a method of treating cancer in a patient are also described. U.S. Pat. No. 6,410,335 discloses a method of predicting risk for prostate cancer, and U.S. Application 2001/0018190 published Aug. 30, 2001 discloses a method of treating prostate cancer with, inter alia, IGFBPs, including IGFBP-3. U.S. Pat. No. 5,840,673 and EP 871,475 disclose a method of inhibiting growth of p53-related tumors by administering IGFBP-3 or a modulator of IGFBP-3 that upregulates IGFBP-3 expression or activity. WO 00/35473 discloses use of IGFBP-6 to treat inflammatory diseases including tumor angiogenesis. WO 94/22466 discloses use of any IGFBP-3 to treat cancer, including prostate cancer. WO 92/14834 discloses two IGFBPs isolated from rat serum, one identified as IGFBP-5 useful in treating cancer.\nExploitation of the interaction between IGF and IGFBP in screening, preventing, or treating disease has been limited, however, because of a lack of specific antagonists. The application of an IGF-1/IGF-2 antagonist as a potential therapeutic adjunct in the treatment of cancer is described by Pietrzkowski et al., Cancer Res., 52: 6447–6451 (1992). In that report, a peptide corresponding to the D-region of IGF-1 was synthesized for use as an IGF-1/2 antagonist. This peptide exhibited questionable inhibitory activity against IGF-1. The basis for the observed inhibition is unclear, as the D-region does not play a significant role in IGF-1 receptor (IGF-1R) binding but rather, in IGF-1 binding to the insulin receptor (Cooke et al., Biochem., 30: 5484–5491 (1991); Bayne et al., J. Biol. Chem., 264: 11004–11008 (1988); Yee et al., Cell Growth and Different., 5: 73–77 (1994)). IGF antagonists whose mechanism of action is via blockade of interactions at the IGF-1 receptor interface may also significantly alter insulin action at the insulin receptor, a disadvantage of such antagonists.\nCertain IGF-1 antagonists have also been described by WO 00/23469, which discloses the portions of IGFBP and IGF peptides that account for IGF-1 and IGFBP binding, i.e., an isolated IGF binding domain of an IGFBP or modification thereof that binds IGF with at least about the same binding affinity as the full-length IGFBP. The patent publication also discloses an IGF antagonist that reduces binding of IGF to an IGF receptor, and/or binds to a binding domain of IGFBP.\nAdditionally, EP 639981 discloses pharmaceutical compositions comprising short peptides that function as IGF-1 receptor antagonists. The peptides used in the pharmaceutical compositions consist of less than 25 amino acids, comprise at least a portion of the C or D region from IGF-1, and inhibit IGF-1-induced autophosphorylation of IGF-1 receptors. Methods of inhibiting cell proliferation and of treating individuals suspected of suffering from or susceptible to diseases associated with undesirable cell proliferation such as cancer, restenosis, and asthma are disclosed.\nGeneration of specific IGF-1 antagonists has been restricted, at least in part, because of difficulties in studying the structure of IGF and IGFBP. Due to the inability to obtain crystals of IGF-1 suitable for diffraction studies, for example, an extrapolation of IGF-1 structure based on the crystal structure of porcine insulin was the most important structural road map for IGF-1 available (Blundell et al., Proc. Natl. Acad. Sci. USA, 75:180–184 (1978)). See also Blundell et al., Fed Proc., 42: 2592 (1983), which discloses tertiary structures, receptor binding, and antigenicity of IGFs. Based on studies of chemically modified and mutated IGF-1, a number of common residues between IGF-1 and insulin have been identified as being part of the IGF-1R-insulin receptor contact site, in particular the aromatic residues at positions 23–25. Using NMR and restrained molecular dynamics, the solution structure of IGF-1 was recently reported (Cooke et al., supra). The resulting minimized structure was shown to better fit the experimental findings on modified IGF-1, as well as the extrapolations made from the structure-activity studies of insulin. Further, De Wolf et al., Protein Sci., 5: 2193 (1996) discloses the solution structure of a mini-IGF-1. Sato et al., Int. J. Pept., 41: 433 (1993) discloses the three-dimensional structure of IGF-1 determined by 1 H-NMR and distance geometry. Torres et al., J. Mol. Biol. 248: 385 (1995) discloses the solution structure of human IGF-2 and its relationship to receptor and binding protein interactions. Laajoki et al., J. Biol. Chem., 275: 10009 (2000) discloses the solution structure and backbone dynamics of long-(Arg(3)) IGF-1.\nPeptide sequences capable of binding to insulin and/or insulin-like growth factor receptors with either agonist or antagonist activity and identified from various peptide libraries are described in WO 01/72771 published Oct. 4, 2001.\nU.S. Application No. 2003/0092631 published May 15, 2003 discloses peptides that antagonize the interaction of IGF-1 with its binding proteins, insulin receptor, and IGF receptor. These IGF antagonist peptides are useful in treating disorders involving IGF-1 as a causative agent, such as, for example, various cancers.\nRegarding the structural information on the classical IGFBPs, they have a molecular mass ranging from 22 to 31 kDa and contain a total of 16–20 cysteines in their conserved amino-and carboxy-terminal domains (Bach and Rechler, supra; Clemmons, Cytokine Growth Factor Rev., 8: 45–62 (1997); Martin and Baxter, Curr. Op. Endocrinol. Diab., 16–21 (1994)). The central domain connecting both cysteine-rich regions is only weakly conserved and contains the cleavage sites for IGFBP-specific proteases (Chernausek et al., J. Biol. Chem., 270: 11377–11382 (1995); Clemmons, supra; Conover, Prog. Growth Factor Res., 6: 301–309 (1995)). Further regulation of the IGFBPs may be achieved by phosphorylation and glycosylation (Bach and Rechler supra; Clemmons, supra). There is no high-resolution structure available for any intact member of the IGFBP family. U.S. Pat. No. 6,500,635 discloses IGFBP-5 and its variants. U.S. Pat. Nos. 6,391,588 and 6,489,294 disclose truncated C-terminal IGFBP-5 fragments with reduced affinity for IGF-I as compared to full-length IGFBP-5. U.S. Pat. No. 6,369,029 discloses stimulating osteogenesis with C-terminal-truncated IGFBP-5. These compounds may be used for stimulating bone cell growth, for treating a bone disorder, or for stimulating mitogenic activity.\nThe NMR structures of two N-terminal fragments from IGFBP-5 that retain IGF-binding activity have recently been reported, showing that residues 40–92 of IGFBP-5 comprise the IGF binding site in the N-terminal domain of that protein (Kalus et al., EMBO J. 17: 6558–6572 (1998)). Other studies have found that N-terminal fragments (residues 1–88 and 1–97) of IGFBP-3 are also able to bind IGFs (Galanis et al., Journal of Endocrinology, 169 (1): 123–133 (2001); Vorwerk et al., Endocrinology, 143 (5): 1677–1685 (2002)).\nIn particular, Galanis et al. have synthesized both the amino-terminal (residues 1–88; N-88) and carboxyl-terminal (residues 165–264; C-165) domains of human IGFBP-3 in bacteria, as fusion proteins with a carboxyl-terminal FLAG peptide. Although only C-165 showed binding to IGF-I and -II by solution-binding assays, both N-88 and C-165 demonstrated binding to IGF-I and -II by biosensor analysis albeit with reduced affinities compared with full-length IGFBP-3. Only the carboxyl-terminal fragment (C-165) was able to form hetero-trimeric complexes with IGF-I and the acid-labile subunit (ALS).\nVorwerk et al. measured the binding of IGF-I and IGF-II to recombinant human N-terminal (residues 1–97; N-97) and C-terminal (residues 98–264; C-98) IGFBP-3 fragments and compared it with IGF binding to intact IGFBP-3 using biosensor analysis. Experiments were carried out either with binding protein or fragment immobilized or with IGF immobilized. These experiments showed that IGF-I and IGF-II bind to IGFBP-3 with affinities of 4–5×10−9 M and similar binding kinetics. The affinities of both N-97 and C-98 for IGF proteins were approximately three orders of magnitude less than that of full-length IGFBP-3.\nU.S. Application No. 2003/0161829A1 published Aug. 28, 2003 and WO 03/025121 disclose fragments of IGFBP-3 that do not bind IGF-I to treat conditions characterized by immune stimulation rather than deficiency. The peptides target the CD74-homology domain sequence at the C-terminus of IGFBP-3 and activities localized to that region, having unique antigenicity. Peptides made to sequences in this region have previously been shown to interfere with the binding of IGFBP-3 to a number of its known ligands, including RXR-alpha, transferrin, ALS, plasminogen, fibrinogen and pre-kallikrein (Liu et al., J. Biol. Chem., 275: 33607–33613(2000); Weinzimeret al., J. Clin. Endocrinol. Metab., 86:1806–13 (2001); Campbell et al, Am. J. Physiol., 275: E321–E231 (1998); Campbell et al., J. Biol. Chem., 274: 30215–30221 (1999); Firth, et al., J. Biol. Chem., 273: 2631–2638, (1998)).\nThe IGFBP-3-derived met al-binding domain peptides differ from previously disclosed IGFBP-3-derived molecules including in their inability to bind IGF-I, their unique antigenicity, and the absence of the IGFBP-3 putative death receptor (P4.33) interaction domain of IGFBP-3 (so-called “mid-region”; amino acids 88–148). The P4.33 putative death receptor is described in WO 01/87238 (Genbank Accession Number BC031217; gi:21411477). For example, WO 02/34916 teaches the use of point mutants of IGFBP-3 in which the binding to IGF-I is impaired. However, the described molecules contain the mid-region of IGFBP-3 and would be expected to exert biological effects by interacting with the P4.33 putative receptor. WO 01/87238 teaches the use of P4.33 modulators for treating disease. The metal-binding peptides do not include the P4.33 putative interaction domain (mid-region of IGFBP-3).\nU.S. Pat. No. 6,417,330, WO 99/63086, and U.S. application No. 2002/0072589 disclose IGFBP-3 variants modified to be resistant to hydrolysis. Also disclosed are variant IGFBP-3s where the nuclear localization signal (NLS) in native IGFBP-3 is altered. Additionally, amino-terminally extended IGFBP-3s are disclosed that include a variety of N-terminal extensions, including peptide and nucleotide binding domains, specific binding members such as ligand-binding domains from receptors or antigen-binding domains from immunoglobins, and peptide and protein hormones and growth factors. N-terminally extended IGFBP-3s may comprise hydrolysis-resistant or NLS-variant IGFBP-3s.\nSome recent publications have described the use of IGFBP-3 peptides to treat cells in culture. The only peptides found to be active on breast cancer cells are derived from the mid-region of IGFBP-3 (McCaig et al., Br. J. Cancer, 86: 1963–1969 (2002); Perks et al., Biochim. Biophys. Res. Comm. 294: 988–994 (2002)).\nU.S. Application No. 2003/0059430 published Mar. 27, 2003 discloses that the IGF-binding protein-derived peptides described above, including short peptides containing just 12–22 amino acids from the C-terminal domain of IGFBP-3, can mimic the full molecule's co-apoptotic, cell-penetrating, and metal-binding properties.\nWO 03/052079 discloses mutants of IGFBP-3 that can inhibit DNA synthesis, can induce apoptosis, bind to neither human IGF-I nor human IGF-II, and comprise a mutation at Y57.\nWO 02/098914 discloses a crystal suitable for X-ray diffraction, comprising a complex of IGF-I or -II and a polypeptide consisting of the amino acids 39–91 of IGFBP-1, the amino acids 55–107 of IGFBP-2, the amino acids 47–99 of IGFBP-3, the amino acids 39–91 of IGFBP4, the amino acids 40–92 of IGFBP-5, or the amino acids 40–92 of IGFBP-6 or a fragment thereof consisting at least of the 9th to 12th cysteine of IGFBP-1, IGFBP-2, IGFBP-3, IGFBP-4, or IGFBP-5 or at least of the 7th to 10th cysteine of IGFBP-6; methods for the determination of the atomic coordinates of such a crystal; IGFBP mutants with enhanced binding affinity for IGF-I and/or IGF-II, and methods to identify and optimize small molecules that displace IGFs from their binding proteins.\nWO 02/34916 discloses mutant IGFBP-3 polypeptides and fragments thereof that have either no binding, or diminished binding to IGFs, yet retain their ability to bind to the human IGFBP-3 receptor (“P4.33”). The fragments are N-deletion fragments with 87 to 264 amino acids. The fragment 1–87 binds poorly to IGF-I, and the other fragments (1–46, 1–75, and 1–80) do not bind at all.\nWO 00/23469 discloses IGFBP fragments that account for IGF-IGFBP binding. It provides an isolated IGF binding domain of an IGFBP or modifications thereof, which binds IGF with at least about the same binding affinity as the full-length IGFBP. It also provides an IGF antagonist that reduces binding of IGF to an IGF receptor. It especially relates to IGFBP-2 fragments, but also provides the isolated IGF binding domains of IGFBP-1, IGFBP-3, IGFBP4, IGFBP-5, and IGFBP-6. As with the IGF binding domain of IGFBP-2, the amino acid sequences comprising the IGF binding domain of the other IGFBPs could include modified forms so long as the binding affinity of the binding domain is about the same as that of the comparable native full-length IGFBP.\nWO 99/32620 discloses IGFBP fragments and utilization thereof, i.e., peptides that are characterized in that the amino acid sequence parts thereof correspond to the amino acid sequence of IGFBP. The invention also relates to cyclic, glycosylated, phosphorylated, acetylated, amidated and/or sulfatized derivatives. These include C-terminal domains of IGFBP-3.\nThe use of gene fusions, though not essential, can facilitate the expression of heterologous peptides in E. coli as well as the subsequent purification of those gene products (Harris, in Genetic Engineering, Williamson, R., Ed. (Academic Press, London, Vol. 4, 1983), p. 127; Ljungquist et al., Eur. J. Biochem., 186: 557–561 (1989); and Ljungquist et al., Eur. J. Biochem., 186: 563–569 (1989)). Protein A fusions are often used because the binding of protein A, or more specifically the Z domain of protein A, to IgG provides an “affinity handle” for the purification of the fused protein. It has also been shown that many heterologous proteins are degraded when expressed directly in E. coli, but are stable when expressed as fusion proteins (Marston, Biochem J., 240: 1 (1986)).\nFusion proteins can be cleaved using chemicals, such as cyanogen bromide, which cleaves at a methionine, or hydroxylamine, which cleaves between an Asn and Gly residue. Using standard recombinant DNA methodology, the nucleotide base pairs encoding these amino acids may be inserted just prior to the 5′ end of the gene encoding the desired peptide.\nAlternatively, one can employ proteolytic cleavage of fusion proteins (Carter, in Protein Purification: From Molecular Mechanisms to Large-Scale Processes, Ladisch et al., eds. (American Chemical Society Symposium Series No. 427, 1990), Ch 13, pages 181–193)).\nThere is a continuing need in the art for a molecule that can be used to elucidate binding epitopes on IGFBP-3 and other ligands, acts as an IGF antagonist to control the levels of circulating IGF as well as receptor response, for therapeutic or diagnostic uses, and can be used for other therapeutic, diagnostic, or assay purposes."}
{"text": "1. Field of the Invention\nThis invention relates to a novel compound and a chiral smectic C (SC*) liquid crystal composition containing the compound and being useful for liquid crystal display elements.\n2. Description of the Related Art\nAt present, as to liquid crystal display elements, TN (Twisted Nematic) type display mode has been most broadly employed, but as far as the response speed is concerned, such TN type display elements are inferior to emissive type display elements (such as those of electroluminescence, plasma display, etc.). Although various improvements in this respect have been attempted, it appears that improvement to a large extent has not yet been realized. Thus, various liquid crystal display devices based on a different principle from that of TN type display elements have been attempted. As one of such devices, there is a display mode utilizing a ferroelectric liquid crystal (N. A. Clark et al: Applied Phys. lett., 36, 899 (1980)). This mode utilizes the chiral smectic C phase (hereinafter abbreviated to SC* phase) or the chiral smectic H phase (hereinafter abbreviated to SH* phase) of the ferroelectric liquid crystal, and those having these phases in the vicinity of room temperature are preferred.\nThese chiral smectic liquid crystal materials may be obtained by blending a plurality of single compounds each exhibiting a chiral smectic phase by itself, but it is known that the materials may be also obtained by adding an optically active liquid crystal compound, preferably a chiral smectic liquid crystal compound to an achiral smectic liquid crystal (exhibiting smectic C phase (SC phase), smectic H phase (SH phase), etc.).\nVarious kinds of compounds exhibiting SC phase have been known, but as to whether or not chiral smectic liquid crystal materials obtained by adding an optically active liquid crystal to the above-mentioned known compounds exhibit superior performances in liquid crystal display utilizing ferroelectricity, no ultimate evaluation thereof has been yet obtained. This is due to the fact that liquid crystal display utilizing ferroelectricity has not yet been technically completed. Thus it is necessary in the present situation to test various novel SC materials."}
{"text": "The disclosure of Japanese Patent Application No. HEI 11-87167 filed on Mar. 29, 1999 including the specification, drawings and abstract are incorporated herein by reference in its entirety.\n1. Field of Invention\nThe invention relates to a control system for starting an engine installed in a vehicle, and more particularly to a control system provided with a plurality of starting systems.\n2. Description of Related Art\nAn engine, such as a gasoline engine is constructed such that fuel-air mixture goes through an intake stroke and a compression stroke before being ignited to cause an explosive combustion, and the pressure generated thereby moves the piston to generate power, and the power is then transmitted to a crankshaft. The control system for the engine is provided with a starting system for rotating the crankshaft at a rotational speed equal to or greater than a predetermined rotational speed when starting the engine.\nJapanese Patent Laid-Open No. HEI 9-117012 discloses a vehicle provided with such a starting system. The vehicle includes an electric motor, an engine, and a generator, which serve as power sources for driving the vehicle, and the generator is connected to the engine. The generator is operated as an electric motor to rotate the engine from a stopped state up to a predetermined idle speed. When the engine speed reaches the idle speed, the fuel is supplied to the engine and then ignited. Then, the generator is no longer operated as an electric motor.\nIn addition to the aforementioned publication, there is also suggested a start control system of a vehicle equipped with an electric motor for driving auxiliaries of a vehicle and a starter motor provided exclusively for engine start, wherein at least one of the electric motor and the starter motor can be used for starting the engine.\nHowever, because a single generator is used to serve as the electric motor for starting the engine, the startability of the engine is adversely affected due to nonconformance of the rotational speed or the torque of the electric motor with the state of the vehicle. Also, because a single generator is used to serve as the electric motor, the startability of the engine is adversely affected when the state of charge of the battery for supplying electrical energy to the electric motor is insufficient to start the engine. Moreover, a pair of large and small spur gears are usually employed as means for power transmission between the starter motor and the engine. Accordingly, there exists the possibility of noise being generated at engine start depending on specifications of the spur gear mesh and the selection of the engine starting mode.\nThe invention has been made in the light of the aforementioned circumstances and therefore it is an object of the invention to provide a start control system capable of improving the engine startability in conformance with the state of a vehicle.\nIn a first aspect of the invention, an engine start control system includes a plurality of starting systems capable of starting the engine, and a controller. The controller determines a first associated amount associated with the rotational resistance of the engine, and then selects the mode of starting the engine by the plurality of starting systems in accordance with the determination results. The first associated amount refers to a phenomenon or a state which allows a direct or indirect determination of the engine rotational resistance.\nAccording to the first aspect, the first associated amount associated with the engine rotational resistance is determined, and the starting mode of the engine is selected based on the determination result of the first associated amount. For example, the starting mode of the engine is selected such that a torque corresponding to the rotational resistance of the engine is obtained. That is, when the rotational resistance of the engine is comparatively high, the starting mode of the plurality of starting systems is selected so that a high starting torque can be obtained. This improves the startability of the engine irrespective of the rotational resistance of the engine.\nThe first associated amount may be the viscosity of the engine oil or the temperature of the engine cooling water. The rotational resistance of the engine is estimated based on the viscosity of the engine oil or the temperature of the engine cooling water.\nIn a second aspect of the invention, the engine start control system provided with a plurality of starting systems capable of starting the engine is also capable of driving a functional device other than the engine using either one or more of the starting systems. The controller of the engine start control system further determines a drive state of the functional device and selects the mode of starting the engine by the plurality of starting systems based on the determined drive state. The drive state of the functional device may be the rotational speed of the functional device, for example.\nAccording to the second aspect, the drive state of the functional device is determined, and the mode of starting the engine by the plurality of starting systems is selected based on the determination result. For example, the starting mode of the engine is selected so that a rotational speed appropriate for starting the engine can be obtained. That is, when the functional device is driven by either one or more of the starting systems and when the rotational speed of the starting system is comparatively high, the engine can be rotated by a starting system that is driving a functional device after the engine is rotated by a starting system that is not driving a functional device. This makes it possible to gradually increase the engine rotational speed, whereby the startability of the engine is improved.\nAccording to a third aspect of the invention, an engine start control system includes a plurality of starting systems capable of starting the engine, a plurality of energy sources which supply energy to respective starting systems, and a controller that determines the states of the plurality of energy sources, and selects the mode of starting the engine by the plurality of starting systems based on the determined states of the plural energy sources.\nAccording to the third aspect, the states of the plurality of energy sources are determined, and the mode of starting the engine by the plurality of starting systems is selected based on the determination result. For example, the starting mode of the engine is selected so that an energy appropriate for starting the engine can be obtained. That is, by selecting a starting mode for starting the engine with a starting system having a higher energy state, the startability of the engine can be improved and a decrease in the energy amount of a starting system having a lower energy state can be suppressed.\nAccording to a fourth aspect of the invention, an engine start control system includes a power transmission member which transmits power from the engine or from at least one of the power sources other than the engine, a plurality of starting systems capable of starting the engine, and a controller that determines a second associated amount associated with the vehicle speed in a state where the vehicle is run by power from the power source, and then selects for selecting the mode of starting the engine by the plurality starting systems based on the determined second associated amount. The second associated amount refers to a phenomenon or a state which allows a direct or indirect determination of the vehicle speed.\nAccording to the fourth aspect, a difference between the engine rotational speed and the rotational speed of the power transmission member is estimated based on the second associated amount associated with the vehicle speed, and the starting mode of the engine is selected based on the estimation result.\nIn accordance with the first through fourth aspects of the invention, the selection of the starting mode of the engine involves selection of which starting system function (rotational speed, torque, and the like) to use, the number of starting systems employed for starting the engine, and the timings for driving and stopping each starting system when a plurality of starting systems are used.\nAccording to a fifth aspect of the invention, an engine start control system includes a first starting system capable of starting the engine by transmitting power to the engine via a gear transmission, a second starting system capable of starting the engine by transmitting power from the engine via a belt transmission, and a controller that, upon starting the engine, first drives the first starting system, then drives the second starting system and continues driving the second starting system even after the first starting system is stopped.\nAccording to the fifth aspect, upon engine start, first, the first starting system is driven and the power is transmitted to the engine via the gear transmission. Then, the second starting system is driven, and the power is transmitted to the engine via the belt transmission. The engine is driven by the power transmitted from the second starting system after the first starting system is stopped. For example, the drive time of the first starting system, which transmits power to the engine via the gear transmission, is set short, so that the generation of noise due to engine start can be suppressed. The belt transmission, which transmits power from the second starting system to the engine, prevents generation of gear noise and abnormal sounds due to backlash as compared to the gear transmission upon engine start. Also, the belt transmission does not require lubricating oil for cooling and lubricating the engaging portions because the device does not contain any engaging portions such as in the gear transmission. Accordingly, temperature increase or power loss caused by agitation of the lubricating oil is eliminated. Also, the construction of the belt transmission is simple because it does not require a lubricating system.\nAccording to a sixth aspect of the invention, an engine start control system includes a first starting system connected to a crankshaft of the engine, a second starting device connected to the crankshaft of the engine, and a controller that, upon starting the engine, first drives the first starting system, then drives the second starting system before driving of the first starting system is terminated.\nAccording to the sixth aspect, the engine is started by drive of the first and the second starting systems in parallel, whereby the startability of the engine is improved."}
{"text": "1. Technical Field\nThis application relates to surgical staplers, and more particularly, to an apparatus for sequentially applying a plurality of surgical fasteners to body tissue.\n2. Background of Related Art\nSurgical devices wherein tissue is first grasped or clamped between opposing jaw structure and then joined by means of surgical fasteners are well known in the art. In some instruments a knife is provided to cut the tissue which has been joined by the fasteners. The fasteners are typically in the form of surgical staples however, two part polymeric fasteners are also utilized.\nInstruments for this purpose can comprise two elongated members which are respectively used to capture or clamp tissue. Typically, one of the members carries a cartridge which houses a plurality of staples arranged in at least two lateral rows while the other member comprises an anvil which defines a surface for forming the staple legs as the fasteners are driven from the cartridge. Where two part fasteners are used, this member carries the mating part, e.g. the receiver, to the fasteners driven from the cartridge. Generally, the stapling operation is effected by a pusher which travels longitudinally through the cartridge carrying member, with the pusher acting upon the staples to sequentially eject them from the cartridge. A knife may travel with the pusher between the staple rows to longitudinally cut and/or open the stapled tissue between the rows of staples. Such instruments are disclosed in U.S. Pat. No. 3,079,606 and U.S. Pat. No. 3,490,675.\nA later stapler disclosed in U.S. Pat. No. 3,499,591 applies a double row of staples on each side of the incision. This is accomplished by providing a cartridge assembly in which a cam member moves through an elongate guide path between two sets of staggered staple carrying grooves. Staple drive members are located within the grooves and are positioned in such a manner so as to be contacted by the longitudinally moving cam to effect ejection of the staples. Other examples of staplers are disclosed in U.S. Pat. Nos. 4,429,695, 5,065,929, and 5,156,614.\nMany of the prior art linear stapling devices discussed above include a significant number of moving parts, small components and machined structural elements that are costly to fabricate and time consuming to assemble. These factors add to the overall cost of the stapling devices and thus increase the costs incurred by hospitals and health care professionals, and ultimately, the patient upon which the devices are utilized. A linear stapler that could be manufactured and assembled less expensively would provide great benefits.\nThe subject application is directed to a linear surgical stapler that is constructed from fewer, less expensive components than known prior art staplers configured to sequentially apply a plurality of surgical staples to body tissue. The stapler includes a first body portion supporting an anvil plate which defines a fastener forming surface and a second body portion configured to releasably mate with the first body portion. In the detailed description which follows, the first and second body portions of the stapler are also referred to as the xe2x80x9canvil half-sectionxe2x80x9d and xe2x80x9ccartridge half-sectionxe2x80x9d, respectively.\nA disposable loading unit is removably supported in the second body portion and includes a cartridge defining a plurality of slots and a tissue contacting surface, a plurality of surgical fasteners disposed in the slots of the cartridge, a plurality of ejectors or pushers positioned adjacent the surgical fasteners. A wedged actuator provided in the disposable loading unit is positioned and configured to enter and translate through the cartridge to sequentially interact with the pushers.\nThe stapler in a preferred embodiment, further includes an elongated actuation member mounted for longitudinal movement within the second body portion and releasably engagable with the wedged actuator, whereby longitudinal movement of the actuation member causes the wedged actuator to interact with the ejectors, driving the surgical fasteners from the cartridge to be formed against the anvil plate.\nThe first body portion preferably includes an elongate anvil support member and a pivoting lever handle. The anvil plate is preferably formed separate from the anvil support member and includes a plurality of staple forming pockets defining the anvil forming surface. The anvil plate also includes means for engaging the anvil support member during assembly of the surgical stapler to securely fasten the anvil plate to the support member. A notched area is defined adjacent a proximal end of the anvil support member and correspondingly positioned detents are formed adjacent a proximal end of the second body portion. The notched area and the detents cooperate to facilitate relative pivotal movement of the first and second body portions when they are mated with one another.\nPreferably, a pair of upstanding flanges are formed on the disposable loading unit proximal of the tissue contacting surface thereof. The flanges define a structural tissue stop to limit the movement of body tissue. The flanges are also dimensioned to engage a pair of corresponding apertures formed in the anvil plate to maintain the first and second body portions in alignment with one another when the surgical stapler is in a closed or clamped position.\nThe wedged actuator of the disposable loading unit is preferably monolithically formed from a planar piece of sheet metal during a stamping process and includes a planar base and a pair of upstanding parallel cam wedges. An upturned flange is formed at a distal end of the actuation member for releasably engaging a complementary slot formed in the base of the wedged actuator. The wedged actuator preferably further includes an upstanding support flange to which a knife blade is fastened. The knife blade is provided to form an incision in the stapled body tissue.\nA retaining channel depends from a distal end of the second body portion for supporting the disposable loading unit. Preferably, the disposable loading unit and the retaining channel include complementary engagement structures for releasably securing the disposable loading unit in the retaining channel. Opposed bearing structures are formed in the retaining channel at a proximal end thereof for abutting the anvil support beam when body tissue is clamped between the anvil plate and the tissue contacting surface of the cartridge. The bearing structures serve to inhibit the anvil support beam from bending as a result of the compressive forces generated during clamping.\nFurther features of the surgical apparatus of the subject application will become more readily apparent to those skilled in the art from the following detailed description of the apparatus taken in conjunction with the drawings."}
{"text": "A product picking apparatus picks up product items according to a first spacing (e.g., product items arriving on a conveyor) and releases the product items according to a second spacing (e.g., product items placed into a carton). Known product picking apparatuses have used belts or endless tracks to frictionally engage and propel product pickers in a cyclical movement. The contacting surfaces responsible for generation of the friction that drives the product pickers are typically wear parts. Therefore, their condition is constantly changing. Accordingly, the movement of the product pickers is subject to slight variation over time. Such variations can result in “missed pickups” by one or more product pickers. This can result in production line shutdowns and/or cartons or cases that are “light” on product."}
{"text": "1. Field of the Invention\nThe present invention relates to a lithium battery, and more particularly to a lithium primary battery having coiled electrode members for use in supplying high current.\n2. Description of the Prior Art\nIn a conventional coiled type lithium battery (FIG. 1), an electrode body of a cylindrical shape of the battery is formed in such a manner that a negative electrode plate 1 made of a strip sheet of lithium or lithium alloy and a positive electrode 2 made of metal charge collecting mesh holding a positive electrode agent are laminated through a separator 3 made of a non woven cloth, such as polypropylene fibers, and are coiled in a roll.\nMoreover, in order to avoid collapse of the roll of the electrode members, the outer end part of the separator 3, wound around the outer cylindrical surface of the electrode body, is adapted to contact with the inner cylindrical surface of a negative electrode can except for a portion of an adhesive tape wound around the central portion of the outer cylindrical surface of the roll of the electrode body. At the inner central portion of the roll of the electrode body, namely at the beginning of the winding of the electrode body, a rolled portion 3a of only the separator 3 is formed and the stripe-like positive electrode 2 is disposed outside of the rolled portion 3a and the strip-like negative electrode 1 is disposed outside of positive electrode 2 through the separator 3.\nThe inner end 1a of the strip-like negative electrode plate 1 ends beyond the inner end 2a of the positive electrode plate 2 toward the inner portion of the roll of the electrode body in order to prevent slip of the respective members of the electrode body and to facilitate winding of the electrode body when the electrode body is wound by a winding machine by clamping the inner end of the separator 3 by the rotation spindle of the machine. Accordingly, as shown by 3a, only the separator 3 is rolled along some length (normally 90.degree. to 180.degree.) from the inner end la of the strip-like negative electrode plate 1.\nAs the collecting mesh of the positive electrode, there is used an expanded metal plate made of stainless steel which metal plate is cut by a number of slits directed in one direction and the metal plate is expanded in a direction perpendicular to the slit to open the slits. Conventionally, the stripe-like positive electrode plate is formed in such a manner that the expanded metal is sandwiched by positive electrode agent sheets made of a positive active agent, electrically conductive powder and a binder with the positive electrode agent sheets bonded to the expanded metal by a rolling machine.\nHowever, according to recent developments in the wide field of use of the battery, there are increasing cases of wrong use of the battery of this kind such that the battery, despite the fact that the battery is a primary battery, is erroneously subjected to charging or forced discharging by an external power source. An example of forced discharging, which occurs many times, is that one or more old batteries, which are already used and the discharging capacity thereof is depleted or small, are connected with one or more new batteries of which the discharging capacity is plentiful. In this case, if the battery of which the discharging capacity is depleted is forced to be discharged, the battery is loaded with a minus voltage by the current supplied from the other healthy battery.\nIn the conventional coiled type lithium battery, if such error occurs, the battery is apt to generate heat due to an internal short circuit and there are many cases of the batteries catching fire.\nThe holding of the positive electrode agent in the positive electrode is done mainly by bonding of the positive electrode mix, carried on the expanded metal on both sides of the slit. However, in case of expanded metal, if the slit is opened too wide, the expanded metal is apt to be cut off. Therefore, in general, the opening rate of the expanded metal is small, about less than 70%, whereby the bonding areas are insufficient, the sufficient holding force can not be obtained and therefore there is a problem that the positive electrode mix can be easily separated from the expanded metal at the time of formation of the coiled electrode body.\nIn case the pressing force of the positive electrode agent sheet is increased for increasing the bonding force of the positive electrode agent and electric collecting ability, there is such a problem that rigid deformation of and cut-off in the expanded metal sheet easily occur. There is a further disadvantage that since the edge of the opened slit and the cut off edge are acute in the expanded metal sheet, the separator is apt to be damaged.\nIt has also been discovered that when the lithium batteries are subjected to work using relatively large current, the batteries generate a lot of heat and the temperature of the battery becomes higher at the central portion, and the insulation layer 3a is molten, whereby an internal short circuit occurs at the inner end 1a of the negative electrode plate 1 and the inner end 2a of the positive electrode 2 and a further large amount of current flows there with the generation of further heat, causing the lithium and lithium alloy to become molten resulting in burning or an explosion of the battery."}
{"text": "Organic excreta, whether from domestic, stray or wild animals, often has a substantial adverse impact on otherwise pleasant modem lifestyles. The same also represents a significant health threat, due to the plethora of diseases often found in animal excrement. While nature will, over time, degrade the organic components of excreta, such action is relatively very long-term. Moreover, while the same is occurring, excreta poses a health hazard, gives off unpleasant odors, and appears and is repulsive to most individuals. At the same time, solid excreta also present the possibility of being picked up underfoot and spread, sometimes even into the home, office or commercial environment. Finally, excreta while it is deteriorating naturally over a long period of time tends to damage lawns due to heat and other prolonged effects.\nIn an effort to avoid the dangers and unpleasantness associated with the same, numerous jurisdictions have implemented so-called “pooper scooper” laws which require, for example, a dog owner to clean up after a pet, for example, while it is being given its required daily walk. However, in addition to the fact that pet owners are not all fastidious about cleaning up after their animals, numerous sources of excreta, not controllable by pooper scooper regulations abound. These range from problems created by Canadian geese on golf courses and waterfront properties to the excreta of wild deer beside houses in the suburbs, and leavings of stray animals in urban locations."}
{"text": "Systems and applications for delivering computer-based training (CBT) have existed for many years. However, CBT systems historically have not gained wide acceptance. A problem hindering the reception of CBTs as a means of training workers and learners is the inflexibility of courses and training material to adapt to the need of individual learners.\nSome CBTs were based on hypermedia systems that statically linked content. User guidance was given by annotating the hyperlinks with descriptive information. The trainee could proceed through learning material by traversing the links embedded in the material. The structure associated with the material was very rigid, and the material could not be easily written, edited, or reused to create additional or new learning material.\nNewer methods for intelligent tutoring and CBT systems are based on special domain models that must be defined prior to creation of the course or content. Once a course is created, the course material may not be easily adapted or changed for specific training needs of different learners. As a result, the course often fails to meet the needs of the trainee and/or trainer.\nIn addition, CBT systems do not provide for collaborative learning that is typically associated with classes taught in a classroom. Practically all CBT courses are constructed according to the paradigm of individual learning that at best is enhanced by an environment that enables users to communicate with each other. For example, a typical web-based course may include an elaborate hyper media structure. The structure is embedded into an environment that provides general tools for communication (e.g., chat rooms or video conferencing). However, content and collaboration are not integrated.\nTherefore, for the above and other reasons, new methods and technologies are needed to supplement traditional computer based training and instruction."}
{"text": "1. Field of the Invention\nThis invention relates to a catheter and, more particularly, to an electrode catheter having a remotely curvable tip.\n2. Description of the Related Art\nThe advantages of steerable catheters are known. Such catheters are inserted into a blood vessel or other body area and their distal end guided in various directions to reach areas of the body which would otherwise be inaccessible. Examples of steerable catheters are disclosed in U.S. Pat. Nos. 3,749,086 to Kline et al.; 3,773,034 to Burns et al., 3,552,384 to Pieric and 3,521,620 to Cook.\nA steerable catheter should have both flexibility for steering the catheter tip and rigidity to control steering. Known prior art catheters have been unable to successfully combine these two characteristics, tending instead to sacrifice one of these characteristics at the expense of the other. It is also desirable to limit the direction of bending of the distal end of the catheter to provide better steering control.\nCatheters having electrode means for monitoring parts of the body, such as the heart, by transmitting electrical signals for analysis on an external monitor are also known. Currently a need exists for an electrode catheter which is steerable to access areas of the heart which could not be reached by a non-bendable catheter. It would be desirable to provide an electrode catheter with a steerable tip that combines sufficient flexibility to allow steering of the catheter to desired areas with good torsional rigidity and steering control."}
{"text": "In recent years, optical devices, such as light-emitting diodes (LED), laser diodes, and UV photo-detectors have increasingly been used. Group-III nitride compounds, such as gallium nitride (GaN) and their related alloys, have been known suitable for the formation of the optical devices. The large bandgap and high electron saturation velocity of the group-III nitride compounds also make them excellent candidates for applications in high-temperatures and high-speed power electronics.\nDue to the high equilibrium pressure of nitrogen at typical growth temperatures, it is extremely difficult to obtain GaN bulk crystals. Therefore, GaN layers and the respective LEDs are often formed on other substrates that match the characteristics of GaN. Sapphire (Al2O3) is a commonly used substrate material. FIG. 1 illustrates a cross-sectional view of an LED package component. LED 2, which includes a plurality of GaN-based layers formed on sapphire substrate 4. Sapphire substrate 4 is further mounted on lead frame 6. Electrodes 8 and 10 electrically connect LED 2 onto lead frame 6 through gold wires 12.\nSapphire has a low thermal conductivity. Accordingly, the heat generated by LED 2 cannot be dissipated efficiently through sapphire substrate 4. Instead, the heat is mostly dissipated through the top end of LED 2 and through gold wires 12. However, the necessary length of gold wires 12 to extend to lead frame 6 renders the heat-dissipating efficiency low. In addition, electrode 10 occupies chip area and reduces the total chip area available for LED light output."}
{"text": "Monosaccharide is a generic term for a hydrocarbon having a carbon chain of 3 or more carbon atoms, and in nature there are a number of monosaccharides as much as 40 or more species in a variety of forms including free forms or derivatives of metabolites or biosynthetic products as well as complex carbohydrates, oligo-saccharides or polysaccharides. Industrially most important monosaccharide is D-glucose, of which the process for production has been established on an industrial scale using a method of hydrolysis of starch with an enzyme such as amylase, and has been industrially most widely used in the fields of foods, cosmetics and pharmaceuticals in large quantities. Further, D-fructose which is produced by conversion of D-glucose with isomerase is an industrially important sugar. Other industrially applied process for production of some monosaccharides includes production of L-sorbose from sorbitol by an acetic acid bacterium, production of D-galactose from lactose by β-galactosidase, and production of D-xylose from xylan by acid hydrolysis. A large number of other monosaccharides are difficult to produce in industrially large quantities because of industrial and economical limitation in raw materials, generative reaction and yield, and they almost have not yet been utilized industrially.\nThe present inventors have paid attention to such expensive and rarely available monosaccharides and gave a name of rare sugars to these monosaccharides; and they have continued to assiduously study the generative reaction of rare sugars for a long period of years. The present inventors have conducted a study for development on monomeric rare sugars including both D-psicose and D-allose. They, as mentioned in the non-patent document 1, have developed a process for producing D-psicose directly from D-glucose in combination of an industrially used D-xylose isomerase with D-ketohexose 3-epimerase discovered by the inventors, and thus opened a new route to a process for the continuous production of D-psicose.\nD-Psicose which is generally known to exist in molasses of sweet potatoes in small quantities, is one of industrially expected sugars, particularly in food industry, since it has a sweet taste of good quality similar to D-fructose and is best suited as hardly fermentative and low calorie sugar. That is, D-psicose draws attention as a sugar which does not promote the synthesis of fat and does not accumulate body fat, particularly intraperitoneal fat in comparison with monosaccharides such as D-glucose or D-fructose (non-patent document 2). Further, it has been reported that the available energy value of D-psicose is approximately zero (non-patent document 3).\nThus, based on their success in producing D-psicose, they have succeeded in large scale production of a rare sugar D-allose by isomerization of D-psicose using L-rhamnose isomerase. In this connection, the present inventors have discovered L-rhamnose isomerase (also referred to as “L-RhI”) produced by Pseudomonas stutzeri LL-172 strain, and further elucidated that this enzyme works to catalyze the isomerization reaction of L-rhamnose to L-rhamnurose as well as L-rhamnurose to L-rhamnose. L-Rhamnose isomerase produced by Pseudomonas stutzeri LL-172 (IPOD FERM BP-08593) also acts on isomerization between D-allose and D-psicose, and thus this enzyme is capable of producing D-allose from D-psicose. In this situation, it is necessary to use an enzyme of the Pseudomonas stutzeri LL-172 origin in order to produce D-allose from D-psicose (patent document 1). Because of discovery of this enzyme, practical use of the reaction for production of rare sugars aimed by the present inventors took a great step forward; particularly, this was remarkable in production of D-allose from the substrate D-psicose on a large scale. Since it became apparent that this D-allose has various physiological activities, the research is rapidly developing. At present, it has become apparent that D-allose has an effect such as cancer cell inhibition or anti-oxidant effect important in a medical field (patent document 2).\nThe production of D-allose, however, includes the most expensive bottleneck step for separating D-allose from D-psicose during the course of production. In order to produce D-allose as a single product, the product contained in a mixture of sugar solutions formed due to the property of the enzyme has to be separated from the raw materials. In the reaction with L-rhamnose isomerase, D-psicose and D-allose reach to an equilibrium state of 7:3, from which mixture of sugar solutions D-allose is separated by column chromatography. Though this operation for separation is naturally required as an essential step for purification of a monomeric sugar, it is also true that the time and labor are required as the production scale of rare sugar becomes huge.\nOn the other hand, “sugar solution of grape sugar/fruit sugar” containing D-glucose and D-fructose has practically been used as a mixture of sugars at present. The sugar solution of grape sugar/fruit sugar has industrially produced as alternative sugar having a compositional ratio similar to ordinary sugar by converting approximately a half of D-glucose to D-fructose using xylose isomerase, and is commercialized as gum syrup for use in coffee or tea. The mixed sugar solution of D-glucose and D-fructose is difficult to crystallize and utilized in a liquid form not as crystals. This is considered because of very high solubility of D-fructose in water, requiring much time and labor and producing a disadvantage in production cost; thus, it has been utilized as such without crystallization. The name of isomerized sugar usually indicates a mixed sugar solution of D-glucose and D-fructose. This is owing to the reason that at present only the sugar solution produced by isomerization reaction of D-glucose as raw material with xylose isomerase has been industrially made fit for practical use.    Non-patent document 1: Journal of Fermentation and Bioengineering, 80, p 101, 1995    Non-patent document 2: Asia Pacific J. Clin. Nutr. 10, 233-237,    Non-patent document 3: J. Nutr. Sci. Vitaminol 48, 77-80, 2002    Patent document 1: PCT WO2004/063369 A1)    Patent document 2: PCT WO03/097820 A1"}
{"text": "The present disclosure relates to crystalline silicon solar cell manufacturing, and more particularly, to a method for texturing a crystalline silicon substrate utilizing an aqueous alkaline etchant solution which includes a nanoparticle slurry as one of the components thereof.\nA photovoltaic device is a device that converts the energy of incident photons to electromotive force (e.m.f.). Typical photovoltaic devices include solar cells, which are configured to convert the energy in the electromagnetic radiation from the Sun to electric energy.\nIn a typical solar cell, single-crystalline silicon is generally used as one of the components of the cell. In such applications, the single-crystalline silicon needs to have a non-planar surface to improve light capture. Typically, the non-planar surface has concave and convex patterns with a minute pyramid (i.e., square pyramid) shape. In such solar cells, the light reflected from one spot impinges again to another spot on the surface of the crystalline solar cell by virtue of the ‘textured’ surface, penetrating into the solar cell to be effectively absorbed in the solar cell. Although a portion of the impinging light that has not been fully absorbed, but arrives at the back face of the solar cell, is reflected back to the surface again, that portion of impinging light can be reflected again at the surface comprising steeply inclined pyramidal surfaces, thereby confining the light in the solar cell to improve absorption of light and to enhance power generation.\nIn conventional single-crystalline silicon solar cells, the textured structure is formed by immersing the exposed (100) face of a single-crystalline silicon wafer into an alkaline solution, i.e., sodium hydroxide (NaOH) or potassium hydroxide (KOH), which may also include 5 to 30% by volume of isopropyl alcohol and/or 0.3 to 1.5% by weight of silicon. The silicon is etched into the bath as a bath conditioner.\nThe etching rate in anisotropic etchants of the kind described above depends on the crystallographic orientation of the silicon surface being etched. The etching rate on the (111) face is significantly lower than the other crystallographic orientations. Accordingly, the (111) face with the slowest etching rate is advantageously left on the surface. Since this (111) face is inclined by about 54 degree against the (100) face, pyramidal projections constituted of the (111) face and its equivalent faces are formed. The pyramid size and density depends on the KOH or NaOH concentration, the amount of silicon already dissolved in the bath, and additive such as isopropyl alcohol."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an apparatus and method for the testing of a USB device, and in particular, to a USB device whose USB communication ports can also be used for testing of the USB device.\n2. Background Art\nUSB is commonly used to interface data communications or connections between a host or hub (such as a personal computer (PC)) and peripheral devices (such as printers and scanners, among others). FIG. 1 is a simple schematic diagram illustrating a conventional USB device 10, such as a scanner, and how the USB device 10 is tested by the manufacturer of the USB device 10 prior to delivery of the USB device 10 to a customer.\nThe USB device 10 has an integrated circuit (IC) 12. Any conventional IC can be utilized, but in many cases, the IC can be a conventional ASIC as is well-known in the art. The IC 12 is coupled to a motor 14 to control the operation of the motor 14. In addition, the IC 12 is coupled to a sensor 16, which can be CCD sensor. The sensor 16 is utilized by the scanner 10 to scan documents, as is well-known in the art, and the signals representing the scanned information are provided to the IC 12 for processing.\nAlthough the USB device 10 is being described herein as being a scanner, the principles of the present invention are equally applicable to any device or peripheral which has the capability for USB communication with a host device. Such devices and peripherals include keyboards, computer mice, printers, monitors, etc. In most cases, the host device would be a PC, which operates in conjunction with an operating system such as xe2x80x9cWindowsxe2x80x9d(trademark) to issue a series of commands to the USB device 10 to control the operation of the USB device 10. As a result, each USB device 10 requires one corresponding PC to enable the USB device 10 to be operated.\nThe testing of a USB device 10 usually includes a conventional burn-in test as is well-known in the art. FIG. 1 illustrates a testing unit 20 that has a processor 22 that is coupled via conventional USB lines (D+, Dxe2x88x92) to the IC 12 to perform the testing. The processor 22 is typically a PC operating in a xe2x80x9cWindowsxe2x80x9d(trademark) operating system (such as Windows 98). However, since each USB device 10 needs to be separately controlled by a dedicated PC, the manufacturer of the USB device 10 will need to have on hand a large number of PCs to perform testing. For example, a scanner manufacturer may need to utilize up to a few hundred PCs. Since each burn-in test can take up to three hours to complete for each USB device 10, the testing becomes cumbersome, time-consuming, and therefore quite expensive. As a result, many USB devices 10 are not tested prior to release.\nTo overcome this problem, some manufacturers have attempted to use smaller and less complex processors instead of an entire PC system to perform the testing. These smaller processors are less expensive than a conventional PC system because they omit many of the components (e.g., certain boards, cards, etc.) of a PC, and do not need to be supported by a xe2x80x9cWindowsxe2x80x9d(trademark) operating system. As a result, the manufacturer may be able to afford to hook up a few hundred smaller processors at any given time to perform the needed testing for the USB devices 10. One example of such a smaller processor is an MCU 8051, manufactured and sold by Intel Corporation, which is well-known in the art. To use the 8051, the USB device 10 would have two dedicated testing ports (such as an 12C protocol bus) that are only used for testing purposes. These testing ports are usually pins and connectors that are typically provided on the IC 12, and the processor would issue commands to the IC 12 via these testing ports to perform the testing procedures. Unfortunately, these dedicated testing ports cannot be used during the normal use and operation of the USB device 10. These open ports are often aesthetically unpleasant, and are often not tolerated by consumers. In addition, these open testing ports use up precious space on an IC 12, and increase the complexity of the design of the IC 12.\nThus, there still remains a need for an apparatus and method for testing USB devices that avoid the above-mentioned problems.\nIt is an object of the present invention to provide an IC and method for use in a USB device that does not need to be tested by PCs.\nIt is another object of the present invention to provide an IC and method for use in a USB device that does not utilize dedicated testing ports.\nIt is yet another object of the present invention to provide an IC and method for use in a USB device where the same pins and connectors can be used for both testing and normal operation.\nTo accomplish the objectives of the present invention, there is provided an apparatus and method of testing a USB device. The present invention provides a first USB communication line and a second USB communication line. Separate first and second paths are created between an external device and a controller in the USB device, with both paths including the first and second USB communication lines. The apparatus detects whether a USB connection exists over the first and second USB communication lines. If a USB connection exists over the first and second USB communication lines, USB communication is routed via the first path. Otherwise, testing communication is routed via the second path. The external device can be either a USB host or a testing unit."}
{"text": "1. Field of the Invention\nThis invention relates to the measurement of the flow rate of weak cell liquor produced in a chlor-alkali cell. It further concerns the use of a weir-type flow meter for the purpose indicated.\n2. Description of the Prior Art\nIn the operation of a cell room containing a plurality of chlor-alkali cells, it has not been usual to determine, regularly and with accuracy, the rate of production of weak cell liquor from an individual cell. When such a determination has been made, it has been usual to use merely a bucket and stop watch. Such a practice is inconvenient and somewhat dangerous, in that it involves the handling of a hot (approximately 90.degree. Centigrade) solution of caustic.\nThe rate of production of weak cell liquor from an individual cell gives an important indication of how well or poorly an individual cell is currently performing. In one kind of commercial cells of the diaphragm type, the flow rate of weak cell liquor being withdrawn from the cell may range from about two to seven liters per minute. In other cells, the range may be somewhat different. Typically, a cell which has recently been renewed by having its diaphragm replaced will give a flow rate on the higher end of the range mentioned above, and as the cell operates for several weeks or months, the observed flow rate tends to decrease. From time to time, it is possible, in the operation of a diaphragm-type cell, to adjust (increase slightly) the level of the head on the anolyte side of the cell and thereby increase the flow rate. Ultimately, however, because of the buildup of a gel or of insoluble particles of alkaline-earth hydroxides which are formed because of alkaline-earth metal ions contained as an impurity in the brine fed to the cell, the rate of flow through the diaphragm decreases, even when the maximum available head is used, to a value such that further operation becomes uneconomical. In order to keep a cell room containing a great number of such individual cells operating on a smooth and satisfactory basis, it would be highly desirable to have a convenient means whereby the rate of flow of weak cell liquor leaving each individual cell could be easily and conveniently determined."}
{"text": "1. Field of the Invention\nThe present invention relates in general to the field of information handling system manufacture, use, and distribution, and more particularly to a system and method for configuration of information handling systems as mobile information technology systems.\n2. Description of the Related Art\nAs the value and use of information continues to increase, individuals and businesses seek additional ways to process and store information. One option available to users is information handling systems. An information handling system generally processes, compiles, stores, and/or communicates information or data for business, personal, or other purposes thereby allowing users to take advantage of the value of the information. Because technology and information handling needs and requirements vary between different users or applications, information handling systems may also vary regarding what information is handled, how the information is handled, how much information is processed, stored, or communicated, and how quickly and efficiently the information may be processed, stored, or communicated. The variations in information handling systems allow for information handling systems to be general or configured for a specific user or specific use such as financial transaction processing, airline reservations, enterprise data storage, or global communications. In addition, information handling systems may include a variety of hardware and software components that may be configured to process, store, and communicate information and may include one or more computer systems, data storage systems, and networking systems.\nIncreases in the capabilities of information handling systems and networking technologies has led to an increasing reliance by businesses and other enterprises on information handling systems to perform a variety of tasks. Often, information handling system servers are placed in a building or room that has specialized power and cooling equipment to help ensure a compatible environment for the systems. As examples of specialized power needs, a concentration of information handling systems often consume a considerable amount of power, typically require a steady and reliable power supply and typically have backup power sources in the event that a main power source is interrupted. As examples of specialized cooling needs, a concentration of information handling systems often produces a considerable amount of heat as a byproduct of their operation and typically has dedicated cooling systems to remove the heat. Often, enterprise information technology specialists employ an over-kill approach when purchasing and installing infrastructure for an information handling system data center. For example, cooling systems typically are purchased and installed that have excess cooling capacity to provide a margin of error. Extra cooling capacity also provides infrastructure to accommodate growth in the number of systems that it supports. Similarly, information handling system data center rooms are often built with unused space so that additional information handling systems may be added over time. Typically, as enterprise information handling system needs grow, enterprises add and replace existing systems; thus, maneuvering room is generally needed within a data center room to service existing systems of the data center. Data centers face considerable infrastructure costs if existing cooling, power and space in a data center become inadequate to handle enterprise information processing needs.\nA recent industry trend seeks to simplify the set up and operation of a data center by installing all or most of the elements of the data center in one or more mobile modules, such as shipping containers. Two examples of such systems are the BLACK BOX manufactured by SUN MICROSYSTEMS and the ICE CUBE manufactured by RACKABLE SYSTEMS. Information handling systems are installed in a shipping container along with cooling, power and networking infrastructure. The shipping containers conform to standards designed for shipping freight through intermodal transportation infrastructure, such as ISO standards. For example, a typical shipping container has a length of forty feet, a width of eight feet and a height of nine feet and six inches. Five common ISO standard shipping container lengths are twenty, forty, forty-five, forty-eight and fifty-three feet, although a variety of other sizes may also be used. The shipping container is shipped as normal freight so that, when the shipping container arrives at its destination, the information handling systems are ready to operate within the container. Although shipping container solutions provide the enterprise with improved flexibility in setting up a data center building by essentially providing the building, i.e., the shipping container, with the information handling systems, existing shipping container solutions continue to share the problems faced by conventional data centers with respect to power consumption and cooling. One type of shipping container solution attempts to include cooling, power and networking infrastructure with the information handling systems in a common shipping container. Such solutions face difficulty in manufacture due to the variety of different types of components that are assembled, and difficulty in deployment since a module generally must include cooling resources for a variety of different climates. U.S. Pat. No. 7,278,273 describes cooling modules that are shipping containers having HVAC equipment and that are separate from a shipping container that has information handling systems; however, air movers for moving cold air are included with the information handling system shipping container, with each air mover providing air to an isolated portion of information handling systems. The addition of the air movers with the information handling systems tends to reduce the space available for information handling systems and to increase the complexity of assembly of the shipping module."}
{"text": "The U.S. Department of Health and Human Services has recently launched a new plan to convert the present healthcare information infrastructure to a nationwide, electronic network for healthcare information. Currently, only thirteen percent of hospitals nationwide have electronic systems, and only 14 to 28 percent of other physicians maintain electronic systems.\nElectronic healthcare systems provide many challenges. For example, oftentimes, physicians and hospitals use computer systems for billing, while they use manual filing systems for tracking medical records. Therefore, information can be lost, mislaid, or inputted incorrectly for billing purposes. Recently, different software systems, such as CareRevolution by Electronic Healthcare Systems, have provided single point-of-care products that provide one database for both billing and medical information.\nNevertheless, medical and billing information must still be transferred from the patients to the physicians manually. Further, with patients moving and increased traveling, often the patient may not have all his medical information in one place and/or on hand when he needs medical attention. While current healthcare systems allow information to be maintained and held by the healthcare provider, it is desirable for such information to be portable and held with the patent instead.\nHaving the proper patient medical history has been shown to improve medical care and reduce medical errors. Further, improvements in health information technology have been estimated to produce savings of up to ten percent. Therefore, there is a need for a portable device for healthcare information."}
{"text": "The present invention relates to a rotary pipe and fixed channel auger with intermediate outer supporting elements.\nAs is known there are at present used several industrial open conveying systems of the belt, roller, bucket, truck, chain type, which are affected by noise and difficult maintenance problems, in addition to having a comparatively high installation cost.\nComparatively noiseless conveying systems, which do not cause any alterations of the materials, are the so-called closed systems, including the augers, or Archimedean screws, including a helical surface which turns inside a pipe which may be integral therewith or be separated therefrom.\nThey are conventionally used in several systems for conveying solid materials, either in a granular or liquid condition, or for mixing different materials.\nAlso known is an auger sorter, in which the auger is fixed and has a vertical axis.\nIn this device, the auger operates as a slightly slanted fixed channel, for discharging at a low speed loose materials and the like.\nThe conveying augers usually comprise two different types of assemblies: the first includes the spiral which is separated from the perimetrical pipe and is called, perhaps not properly, auger; the second has its full or partial spiral affixed within the rotary pipe and is usually called \"rotary tube\".\nThe industrial type of auger is mainly used for horizontal conveyance.\nRarely it is also used for conveying materials on slight slopes, these materials consisting of loose solid materials or dense liquids. For conveying granular solid materials there are conventionally used vertical augers.\nThe conveying device is loaded at any set region of the path, and it is likewise unloaded, by opening the bottom of the channel. The drive is provided by a longitudinal central shaft, of hollow construction, usually made of steel and supported by supporting elements spaced from one another by 2.5-3.5 m.\nSeveral types of spirals are used and, more specifically, an iron sheet spiral, a ribbon or belt spiral, a blade spiral including a plurality of blades operating as the blade elements of a propeller, and allowing the advancing speed of the material to be changed by changing the angles of the blades, as well as a fixed-spiral auger made of a manganese steel alloy for conveying materials having high abrading properties.\nOnly by using comparatively great diameter shafts it has been possible to make great length of augers, having a length of up to 50 m; the diameter of the shaft being increased from 140 to 300 mm and that of the outer channel from 400 to 630 mm, the supporting elements being spaced by 7-8 up to a maximum of 10 m and, for greater spacings, several series arranged augers must be provided.\nThese prior auger devices, however, are affected by great problems: in fact, the volumetric mass of material which can be conveyed, that is the conveying capability of the auger is greatly reduced, the power required for driving the auger increasing proportionately to the resistances to the movements and, simultaneously, the device has a poor reliability and requires frequently maintenance operations.\nIndependently from its size, the spiral is interrupted at each supporting elements, with consequent great difficulties from the conveyance standpoint, and with the possibility of jams, as powder or paste materials are conveyed.\nMoreover, the intermediate supporting elements are subjected to a high corrosion and, because of friction, the temperature can excessively increase so as to cause fires or explosions, whereas some types of chemical products tend to cure under high temperatures, so as to cause the above mentioned jams.\nIn order to reduce the maintenance difficulties, prior auger devices include a half-cut bearing element, which can be opened in order to access the support member without disassembling the auger and, usually, several augers are series arranged with consequent greater costs for the motors, electric controls and so on."}
{"text": "Throughout this application various publications are referred to by partial citations within parenthesis. Full citations for these publications may be found at the end of the specification immediately preceding the claims. The disclosures of these publications, in their entireties, are hereby incorporated by reference into this application in order to more fully describe the state of the art to which the invention pertains.\nNeuroregulators comprise a diverse group of natural products that subserve or modulate communication in the nervous system. They include, but are not limited to, neuropeptides, amino acids, biogenic amines, lipids and lipid metabolites, and other metabolic byproducts. Many of these neuroregulator substances interact with specific cell surface receptors which transduce signals from the outside to the inside of the cell. G-protein coupled receptors (GPCRs) represent a major class of cell surface receptors with which many neurotransmitters interact to mediate their effects. GPCRs are characterized by seven membrane-spanning domains and are coupled to their effectors via G-proteins linking receptor activation with intracellular biochemical sequelae such as stimulation of adenylyl cyclase. This application describes the identification of two GPCRs, SNORF62 and SNORF72, as receptors for neuromedin U (NMU) neuropeptides.\nNeuropeptides are synthesized and released from neurons to mediate their effects on cells within the nervous system or on peripheral targets. NMU-25 and NMU-8 are bioactive peptides originally isolated from porcine spinal cord (Minamino, N. et al. 1985a and 1985b). NMU-8 corresponds to the C-terminus of porcine NMU-25 preceded by Arg-Arg residues and may therefore be generated by enzymatic cleavage. NMU homologues have been identified in many species including human (25 amino acids) and rat (23 amino acids).\nThe amino acid sequence for human NMU-25 is as follows:                F R V D E E F Q S P F A S Q S R G Y F L F R P R N—NH2 (SEQ ID NO: 5).        \nThe amino acid sequence for porcine NMU-25 is as follows:                F K V D E E F Q G P I V S Q N R R Y F L F R P R N—NH2 (SEQ ID NO: 6).        \nThe amino acid sequence for rat NMU-23 is as follows:                Y K V N E-Y Q G P-V A P S G G F F L F R P R N—NH2         \n(SEQ ID NO: 7)(—indicates gaps in rat NMU-23 sequence to demonstrate optimum alignment).\nThe amino acid sequence for NMU-8 is Y F L F R P R N—NH2 (SEQ ID NO: 8). All of the preceding sequences were taken from Nandha and Bloom 1993 and Austin, et al. 1995.\nInterestingly, the 8 carboxy-terminal residues of human NMU-25 are identical to those in porcine, rabbit and guinea pig NMU and differ only by one residue from the C-terminus of frog, rat, dog and chicken NMU (Austin et al. 1995). C-terminal NMU peptides (8–9 amino acids) have also been identified in guinea pig, chicken and dog tissue extracts (Minamino et al. 1985a and 1985b, Domin et al. 1989, O'Harte et al. 1991). Indeed, the region of rat NMU-23 critical for smooth muscle contractile activity was found to reside between residues 17–22 (C-terminal region) (Hashimoto et al. 1991, Sakura et al. 1991). However, other groups have demonstrated the necessity of the amidated C-terminal asparagine (Asn23) for activity as well (Nandha and Bloom 1993). Full length NMU is approximately 3-fold more potent than NMU-8 in smooth muscle contraction assays suggesting that the N-terminal region of the peptide also contributes to the activity (Nandha and Bloom, 1993). Several residues in the middle region of the peptide are conserved between species including Glu5, Gln8 and Pro10 supporting the functional importance of this region of the peptide (Nandha and Bloom, 1993). The C-terminus of NMU shares some homology with rat pancreatic polypeptide (PP): Leu-X-Arg-Pro-Arg-X-amide and contains a terminal asparaginamide also present in vasoactive intestinal polypeptide (VIP) (Nandha and Bloom 1993). However, the structure of NMU is unrelated to the other neuromedin peptides isolated by Minamino et al. (1985a,b).\nA profound effect of NMU has been observed in rats on the in vivo release of stress-related modulators from the anterior pituitary and adrenal glands (reviewed in Malendowicz and Markowska 1994). Following a single subcutaneous injection of NMU-8 (6 μg/100 g body weight), adrenocorticotropic hormone (ACTH) blood concentrations are elevated transiently (3–12 hours) and plasma corticosterone levels remain elevated for 24 hours (Malendowicz et al., 1993). In addition, the stress-evoked rise in corticosterone was absent in rats treated for 6 days with NMU-8 (Malendowicz et al. 1994a). Since corticosterone exerts both mineralocorticoid and glucocorticoid effects, regulation of its release by NMU ligands would be expected to modulate fluid homeostasis, ionic balance and metabolism. Although the mechanisms that mediate these effects remain unclear, identification of NMU-like immunoreactivity in nerve fibers in the rat hypothalamic paraventricular and supraoptic nuclei suggest a potential role for NMU in the hypothalamic regulation of pituitary function (Steel et al. 1988).\nThe corticosterone releasing effects of NMU may be mediated in part by direct effects on the adrenal gland. In rat adrenal gland slices, NMU-8 markedly increased basal corticosterone and pregnenolone steroid secretion (Malendowicz et al. 1994a and 1994b). These effects require the presence of adrenal medulla suggesting that NMU-8 acts on medullary chromaffin cells which may stimulate cells of the cortex through a paracrine mechanism. On the other hand, rat NMU-23 directly decreased basal corticosterone secretion from isolated rat inner adrenocortical cells (in the absence of medullary cells) while NMU-8 was without effect (Malendowicz and Nussdorfer 1993). This discrepancy between NMU-23 and NMU-8 effects on adrenal cortical cells suggests that the NMU receptor in these cells differs from that responsible for smooth muscle contraction. Repeated NMU-8 administration also decreased adrenal weight and the number of cells in the zona reticularis, further suggesting a stimulatory role for NMU on adrenal gland (Malendowicz et al. 1994a). NMU ligands may therefore be useful for directly regulating secretion from the adrenal gland.\nAlthough NMU-like immunoreactivity has not been demonstrated within the adrenal gland or circulating in plasma, corticotrophs within the anterior lobe of rat and human pituitary gland contain high levels of NMU-like immunoreactivity (Steel et al. 1988) suggesting a possible hormonal role for NMU. Co-release of NMU with other bioactive peptides is likely to occur since NMU was observed by electron microscopy to be present in the same secretory granules as ACTH and galanin (Cimini et al. 1993).\nFurthermore, both ACTH and NMU are present in human pituitary corticotropinomas as well as in ACTH expressing tumors from a variety of other tissue sources (Steel et al. 1988). Supporting a potential hormonal role of NMU, is the identification of a small population of NMU positive parafollicular C-cells in rat thyroid gland (Domin et al. 1990 and Lo et al. 1992).\nActivities of this peptide also include a hypertensive effect when given intravenously to rats at a high dose (1 nmole; Gardiner et al. 1990). However, at a lower dose (0.1 nmole), NMU caused potent constrictor effects on the superior mesenteric vascular bed reducing mesenteric blood flow without changing systemic blood pressure. The NMU-induced reduction in mesenteric blood flow was also demonstrated in dog (Sumi et al. 1987). In addition, a slight increase in blood flow to the pancreas was measured in these experiments. Such actions suggest the involvement of NMU in the regulation of blood flow to the digestive tract and subsequent effects on digestion.\nNMU was originally isolated based on its potent uterine contractile activity in vitro and has contractile activity on other smooth muscle preparations including chicken crop (Minamino, N. et al. 1985a, 1985b). Isolated muscle strips from the dome of the human urinary bladder were also contracted by NMU (Maggi et al. 1990) suggesting a role for this peptide in urinary control. NMU-like immunoreactivity has been identified in high levels in the rat genito-urinary systems including vas deferens, prostate, fallopian tube, urethra, vagina, ovary and uterus (reviewed in Nandha and Bloom 1993). Smooth muscle contractile or other hormonal effects of NMU in these tissues may regulate urinary control and/or reproductive functions.\nAlong with many other neuropeptides, NMU is present in nerves throughout the gastrointestinal tract (reviewed in Nandha and Bloom 1993). NMU stimulates contraction of isolated longitudinal muscle of human ileum (Maggi et al. 1990) and rat stomach circular muscle (Benito-Orfila et al. 1991) suggesting a role for NMU in gastric emptying and intestinal motility. Interestingly, porcine jejunum (Brown and Quito 1988) and guinea pig small intestine (Minamino et al. 1985b) are not contracted by NMU indicating species differences in gut regulation by this peptide. However, ion transport is modulated by NMU-8 in isolated porcine jejunal mucosa (Brown and Quito 1988). NMU-like immunoreactivity in the intestine has been localized to both the submucosal and myenteric ganglion cells (Ballesta et al. 1988) consistent with the observed effects on contractility, blood flow and absorptive/secretory functions (Ballesta et al. 1988).\nAlthough higher concentrations of NMU are found in the periphery than in the central nervous system (CNS), immunocytochemical analysis demonstrated the presence of NMU in nerve fibers in many CNS regions with concentrations in discrete functional systems (Honzawa et al. 1987, Ballesta et al. 1988 and reviewed in Domin et al. 1987). For example, NMU-like immunoreactivity was identified in all of the cranial nuclei associated with somato-motor function (Ballesta et al. 1988). Several structures associated with sensory processing are also rich in NMU containing fibers including spinal cord (dorsal horn>ventral horn), trigeminal sensory nuclei, vestibular nuclei and other nuclei associated with descending spinal pathways (Honzawa et al. 1987). This localization suggests a role for NMU in perception and processing of sensory stimuli including pain.\nThree cerebellar nuclei (nucleus medialis, interpostitus and lateralis) also demonstrated NMU-like immunoreactivity, consistent with the potential importance of NMU in sensory processing. Neuronal cell bodies containing NMU-like immunoreactivity have been identified in the arcuate nucleus of the hypothalamus, an area identified as important for the regulation of food intake and neuroendocrine control.\nRelatively high levels of NMU-like immunoreactivity were also detected in the nucleus accumbens (Domin et al. 1987), an area where dopaminergic transmission is involved in reward and reinforcement of learned behaviors. The presence of NMU in another area important in dopaminergic transmission, the substantia nigra, (Domin et al. 1987) suggests a role for NMU in the modulation of dopaminergic actions in movement control as well. NMU-like immunoreactivity is also found in the hippocampus, amygdala and other portions of the limbic system suggesting a role for NMU ligands in affective disorders, psychosis and cognition.\nG-protein coupled receptors (GPCR's) activated by this peptide and related analogues were postulated to exist based on binding of [125I]rat NMU in rat uterus membranes (Nandha K. A. et al. 1993). The binding is saturable and of high affinity (Kd=0.35 nM, maximal binding capacity (Bmax)=580 fmol/mg protein). This affinity corresponds to the EC50 of contractile activity in this tissue, 0.2 nM, consistent with the involvement of this binding site in NMU-induced uterine contraction. The GTP analogue, GTPγS, inhibited binding of [125I]rat NMU-23 suggesting that the binding site is a GPCR. In addition, chemical cross-linking identified the binding protein as having an apparent Mr of 48,500 which is consistent with the expected size of a GPCR protein.\nNandha et al. (1994) also identified [125I]rat NMU-23 binding sites in rat uterus tissue slices and in the indusium griseum by autoradiography."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nComposite materials, in particular panels used in aeronautical applications, proximate to the engines, are intended for operation at temperatures below a maximum threshold. These panels, whether forming one single layer or used in a sandwich structure comprising a central core such as honeycomb covered with two external skins or face sheets, include unique mechanical characteristics and a reduced mass, and may take various shapes.\nHowever, subsequent to thermal events comprising an increase of temperature above a maximum threshold, and depending on the exposure times, some damage to the materials of the panels may occur, which may reduce the structural strength characteristics. Thus, it is desirable to be able to characterize the level of the material damage during maintenance operations in order to increase safety.\nIn order to carry out inspection on these composite materials, it is known to apply a primary paint over the surface, in particular a polyurethane paint, which will undergo color transformations during thermal events.\nA first known method for damage assessment consists in visually studying the discoloration of the primary paint in order to deduce therefrom some level of evolution of the material. This method gives a qualitative indication, but does not allow for an accurate quantitative measurement of the damage.\nAnother known non-destructive method, in particular presented by the documents EP-A1-2343531 and US-A1-20110001047, performs an analysis of the discoloration of a primary paint by a technique based on infrared spectrometry.\nIn particular, the second mentioned document presents the following method: after having carried out a calibration model, infrared radiation is emitted on a point of the surface of the material to be analyzed; then, the diffused reflection spectrum of this surface is measured; and by considering on this spectrum a criterion which is compared to the calibration model, a history of the thermal or chemical aspects of the coating is assessed, in order to deduce therefrom a condition of the composite material supporting it.\nNonetheless, this method is complex to carry out, and the obtained results highly depend on the implementation of the primary paint. In particular, thickness variations of the coating layer, of its polymerization cycle, of the number of deposited layers or of the surface condition, will disturb the infrared reflection. The results as well as the conclusions related to the damage to the material are not accurate.\nOther known non-destructive methods for assessing the damage to the materials use in particular ultrasonic or thermographic conventional techniques which allow in particular detecting detachments between the layers of the materials. Thus, it is possible, for example, to control the detachment between the skin and the core of a panel comprising a sandwich structure. Nonetheless, this method does not allow assessing the condition of the material by itself."}
{"text": "Facsimile (or fax) technology is one of the most widely used forms of communication between companies and individuals. Notwithstanding the continuous growth of the popularity of email communications, the fax still holds a strong position as a reliable means of communication. It is presently estimated that more companies in the world have a fax number than an email address.\nAmong the advantages of fax communications are immediate delivery of documents, ability to transmit graphical documents, interoperability of fax devices throughout the world, improved security over email, and easy to use.\nTraditional fax technology utilizes Public Switched Telephone Network (PSTN), also called Global Switched Telephone Network (GSTN), to transfer data between fax devices. Traditional fax devices comply with a variety of protocols, among which are: T.4—“Standardization of Group 3 Facsimile Apparatus for Document Transmission”, ITU-T (CCITT), July 1996; T.6—“Facsimile Coding Schemes and Coding Control Functions for Group 4 Facsimile Apparatus”, ITU-T (CCITT), November 1988; and T.30—“Procedures for Document Facsimile Transmission in the General Switched Telephone Network”, ITU-T (CCITT), July 1996; all are hereby incorporated in their entirety by reference.\nInternet Fax (I-Fax) is a term used to describe extensions of the traditional fax technology that allow the use of the Internet for fax transmission. Examples of I-Fax are fax-over-IP and fax-through-email. Among the standards recommended for use in I-Fax are: T.37—“Procedures for the transfer of facsimile data via store-and-forward on the Internet”, ITU-T (CCITT), June 1998; T.38—“Procedures for real-time Group 3 facsimile communication over IP networks”, ITU-T (CCITT), 1998; F.185—“Internet facsimile: Guidelines for the support of the communication of facsimile documents”, ITU-T (CCITT), 1998; RFC 2301—“File Format for Internet Fax”, IETF, March 1998; RFC 2302—“Tag Image File Format (TIFF)—image/tiff MIME Sub-type Registration”, IETF, March 1998; RFC 3191—“Minimal GSTN address format in Internet Mail”, IETF, October 2001 (obsoletes RFC 2303, March 1998; updates RFC 2846, June 2000); RFC 3192—“Minimal FAX address format in Internet Mail”, IETF, October 2001 (obsoletes RFC 2304, March 1998; updates RFC 2846, June 2000); RFC 3965—“A Simple Mode of Facsimile Using Internet Mail”, IETF, December 2004 (obsoletes RFC 2305, March 1998); RFC 2306—“Tag Image File Format (TIFF)-F Profile for Facsimile”, IETF, March 1998; RFC 2542—“Terminology and Goals for Internet Fax”, IETF, March 1999; and RFC 3297—“Content Negotiation for Messaging Services based on Email”, IETF, July 2002; all are hereby incorporated in their entirety by reference.\nAdvanced capabilities of Internet Fax, such as color fax transmissions, are described, inter alia, in RFC 2531—“Content Feature Schema for Internet Fax”, IETF, March 1999 and RFC 2532—“Extended Facsimile Using Internet Mail”, IETF, March 1999; all are hereby incorporated in their entirety by reference.\nIn a conventional fax-through-email system, a fax service provider receives a fax transmission over the PSTN, converts the received fax transmission to a computer-readable file, e.g. TIFF or PDF (Portable Document Format), and forwards the file to the intended fax recipient via email. If a user sends a fax, the user provides document(s) to be faxed to the fax service provider via email, website, or software. The documents typically are computer-readable files. The fax service provider converts the document(s) to a faxable format (typically 1-bit black and white images), and transmits them to the recipient over the PSTN.\nPresently, the majority of fax devices only support transmission of 1-bit black and white (monochrome) images. Thus, black and white faxes constitute the vast majority of all faxes transmitted. If the user's files to be faxed contain colors or shades of gray (hereinafter, the term color includes shades of gray), the colors will be typically converted to 1-bit black and white for compatibility with the majority of fax devices. At present, colors are either converted to black or dithered. Dithering is a method of simulating colors by producing black dots/pixels to represent brightness of the original color (similar to halftoning in newspaper photographs). Dithering will transform light colors to fewer dots with larger distances between them. The darker colors will be represented by more dots positioned closer to each other (additionally the size of the dots may be increased).\nVarious algorithms have been developed to produce a dithered image that more closely resembles the original image. Among them are error diffusion (dispersion, distribution) algorithms, which dither the image by diffusing the quantization error of a pixel to its neighboring pixels (e.g. the Floyd-Steinberg dithering algorithm).\nVarious colors with the same or similar brightness may look the same after dithering. Referring to FIG. 1, three elements of document A having various colors (represented by different line patterns) may look virtually undistinguishable after dithering in document B. Similarly, a fax document with a colored text on a colored background may become unreadable if the brightness of the text color is similar to the brightness of the background color (FIGS. 6 and 7).\nTherefore, new methods, apparatus, software, and systems are needed to provide document optimization prior to their conversion to black and white. Additionally, new methods, apparatus, software, and systems are needed to provide document optimization for use in fax technology."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for electrically processing image information such as an image file or a facsimile apparatus, and more particularly to such image processing apparatus adapted for processing image information in compressed form, for example image information stored or transmitted in the run-length encoded form using the modified Huffman (M-H) codes.\n2. Description of the Prior Art\nIt has been proposed to read an original image photoelectrically to generate electric image signals, and to transmit such signals to a printer or a file, or to a distant printer apparatus through a transmission line. In the case of such image signal transmission among plural units, satisfactory synchronization is indispensable between the transmitting and receiving units. However, satisfactory image signal transmission cannot be expected by mere synchronization if the signal processing speed is different between the different units.\nAccordingly, there has been considered adjustment the speed or timing of transmission by providing at least either one of the transmitting and receiving units, for example a reader and a printer, with a buffer memory and storing the image signals to be transmitted in such buffer memory.\nSuch a buffer memory requires a certain margin in memory capacity in order to ensure satisfactory synchronization. However an increase in the memory capacity is inevitably limited as it leads to an increase in the cost.\nAlso in the case that a compression process is involved in the image data handling, the quantity of data after compression varies significantly, so that the number of pages that can be stored in a memory fluctuates according to the case. At the transmission of compressed image data, certain other signals such as timing signals are transmitted simultaneously, and the storage of such other signals into the aforementioned memory gives rise to waste of the memory capacity."}
{"text": "1. Field of the Invention\nThis invention relates to an oil pressure booster, and more particularly to an improvement of an oil pressure booster which is used in a fork lift, a truck, an automobile, or the like vehicle.\n2. Description of the Prior Art\nAn oil pressure booster of the prior art, which is used for instance in a fork lift, comprises a throttle valve formed of a spool to be driven by a brake pedal, which throttle valve throttles the passage of operating oil flow being delivered upon depression of the brake pedal for generating an oil pressure upstream thereof, a control valve formed of the aforesaid spool, a return path valve formed of said spool and normally communicating a boost chamber to a drain port therethrough, and a power piston operatively related to said spool. The aforesaid oil pressure is applied to the boost chamber through said control valve for moving the power piston in the direction of depressing the brake pedal so as to follow the movement of the spool. The oil pressure booster of the prior art has a shortcoming in that its operating characteristics with an unoperated power steering connected thereto (open characteristics) is widely different from its operating characteristics at the time when the power steering connected thereto is in operation (closed characteristics). Thus, when the power steering is not operated while the oil pressure booster is operated, if the power steering is put into operation for shifting the open characteristics to the closed characteristics, with a certain depressing force applied to the brake pedal (which relates to a boost chamber pressure), the oil pressure delivered to the boost chamber under the closed characteristics is considerably higher than that under the open characteristics, whereby the pressure difference under the two characteristics tends to push up the brake pedal (to be referred to as \"oil pressure interference\"), which push up cause a discomfort to the driver."}
{"text": "Electrophoresis is a well-known technique for separating macromolecules. In electrophoretic applications, molecules in a sample to be tested are migrated in a medium across which a voltage potential is applied. Oftentimes, the sample is propagated through a gel which acts as a sieving matrix to help retard and separate the individual molecules as they migrate.\nOne application of gel electrophoresis is in DNA sequencing. Prior to electrophoresis analysis, the DNA sample is prepared using well-known methods. The result is a solution of DNA fragments of all possible lengths corresponding to the same total sequential order, with each fragment terminated with a tag label corresponding to the identity of the given terminal base.\nThe separation process employs a capillary tube filled with conductive gel. To introduce the sample, one end of the tube is placed into the DNA reaction vial. After a small amount of sample enters the capillary end, both capillary ends are then placed in separate buffer solutions. A voltage potential is then applied across the capillary tube. The voltage drop causes the DNA sample to migrate from one end of the capillary to the other. Differences in the migration rates of the DNA fragments cause the sample to separate into bands of similar-length fragments. As the bands traverse the capillary tube, the bands are typically read at some point along the capillary tube using one of several detection techniques.\nUsually, multiple DNA preparation reactions are performed in a commercially available microtitre tray having many separate low-volume wells, each holding on the order of 200-1000 micro-liters. The microtitre trays come in standard sizes. In the biotech industry, the currently preferred microtitre tray has a rectangular array comprising of 8 rows and 12 columns of wells. The centers of adjacent wells found in a single row are separated by approximately 0.9 cm, although this figure may vary by one or two tenths of a millimeter. The same holds for the spacing between adjacent wells in a single column. The rectangular array of 96 wells has a footprint within an area less than 7.5 cm.times.11 cm.\nMiniaturization has allowed more wells to be accommodated in a single microtitre tray having the same footprint. New trays having four times the density of wells within the same footprint have already been introduced and are fast becoming the industry standard. Thus, these new trays have 16 rows and 24 columns with an inter-well spacing of approximately 0.45 cm.\nIt is not uncommon to analyze several thousand DNA samples for a given DNA sequencing project. Needless, to say, it is time consuming to employ a single capillary tube for several thousand runs.\nPrior art devices have suggested means for analyzing DNA bands in multiple capillaries simultaneously. Such a device is disclosed in U.S. Pat. No. 5,498,324 to Yeung et al, whose contents are incorporated by reference in their entirety. This reference teaches a means for detecting the DNA bands as they are separated in multiple capillary tubes which are positioned parallel to another. However, in such an arrangement, each capillary tube is filled with gel and a sample is introduced into each capillary tube.\nThe arrangement described above takes a considerable amount of time to fill each capillary tube with gel. It also takes considerable effort to introduce a reaction sample into one end of each of the tubes reproducibly and reliably.\nIt is also not uncommon that one uses the same capillary tube for several consecutive sample runs. This, obviously risks cross-contamination of samples, which is a further disadvantage in certain prior art arrangements."}
{"text": "Certain embodiments of the present invention generally relate to a contact configured to deflect across a defined vertical range while electrically connecting a circuit board to a processor.\nMany large electronic devices, such as computers, use sockets to connect different electronic components. For example, pin grid array (PGA) sockets are used to electrically connect electronic packages, such as processors, to printed circuit boards. PGA sockets facilitate electrical communication between a large number of pins on the processor and contacts on the circuit board. PGA sockets may utilize a cover that is slidably movable on a base between open and closed positions. The sliding movement may be actuated, for example, by a lever. The cover has a hole array configured to match a pin array on the processor. Similarly, the base has an array of pin receiving chambers configured to accept the pin array of the processor. The processor is mated to the socket by first placing the processor such that its pins slide into the holes of the cover. With the cover in the open position, the processor pins pass through the holes of the cover into the pin receiving chambers of the base, but are not electrically connected to the pin receiving chambers of the base.\nWhen the cover is slid to the closed position, the processor pins electrically connect to contacts in the pin receiving chambers in the base. The contacts have fingers that flexible receive the processor pins therebetween. This PGA base and cover arrangement, however, requires use of a mechanism, such as a lever assembly, thereby introducing excess parts and manufacturing cost. The PGA base and cover arrangement also requires additional space for the contacts as the fingers on the contacts must flex outward away from each other to receive the processor pins. These drawbacks are especially troublesome in applications where space is at a premium, such as on motherboards for desktop and laptop computers.\nConsequently, land grid array (LGA) sockets have been developed which require only vertical compression to allow a processor and circuit board to electrically communicate. The LGA sockets do not require the lever mechanism, and can be used in applications with more stringent space requirements. LGA sockets, however, require a vertical compression force to be continuously applied to the processor to maintain proper communication between the processor and the circuit board.\nOne LGA socket has been proposed in an application, entitled “Surface Mount Technology Land Grid Array Socket,” filed on Aug. 5, 2002, and afforded Ser. No. 10/212,414. FIGS. 1 and 2 illustrate the above-noted electrical system 10 including a surface mount land grid array (LGA) socket 11. The electrical system 10 also includes a circuit board 12 to which the socket 11 is mounted and a processor 18 mounted to the socket 11. The socket 11 includes a frame 14, a housing 16 (see FIG. 2), and bias spring arms 20. The frame 14 includes an array of holes 86 therein that hold socket contacts in a pattern that corresponds to a pattern of contacts provided on the bottom of the processor 18. The bias spring arms 20 locate and position the processor 18 with respect to the socket 11 such that the processor contacts align and engage socket contacts to facilitate electrical communication between the processor 18 and the circuit board 12. When the housing 16 is positioned in the frame 14 and the processor 18 is positioned on the housing 16, the processor and socket contacts are placed under a desired vertical load between the circuit board and the processor.\nHowever, existing LGA socket contacts have experienced certain limitations, such as an unduly limited range of deflection. More specifically, after the processor is positioned on top of the socket contact, the processor applies a normal vertical force that deflects the socket contact between first and second contact positions. The range of deflection determines certain tolerances of the individual components, such as the socket, processor, and circuit board. Conventional LGA socket contacts have an unduly limited range of deflection which places undesirably narrow limits on the component tolerances. Additionally, conventional socket contacts may not return to their unbiased first position upon removal of the processor while affording the desired deflection range.\nA need exists for an LGA socket contact that addresses the above noted problems and others experienced heretofore."}
{"text": "1. Field of the Invention\nThe present invention relates to fishing lure storage boxes.\n2. Description of Related Art\nU.S. Pat. No. 5,185,952 to Bruce teaches a box with a display panel having a plurality of tracks, each track holding one or more clips, each clip comprising two resilient bifurcated legs having a hook receiving cavity between them. The hooks of adjacent lures stored in Bruce's box may become entangled unless a precisely positioned clip is provided for each hook.\nIn U.S. Pat. No. 4,958,730 Bunten provides an elastic retainer holding one hook of a multi-hooked lure to the top of a lure drawer. A second hook is retained by a fixed rod adjacent the bottom of the lure drawer.\nIn U.S. Pat. No. 4,240,222 Covington discloses a modular, multi-compartmented tackle box having various means (e.g., a transverse rod, a channel filled with a resilient foam, or a pierceable board) for suspending a lure by one hook. Covington's tackle box, if jostled, may be found to contain entangled lures.\nMalmquist, in U.S. Pat. No. 2,608,459, teaches a tackle box comprising a plurality of separate lure storage compartments, each compartment lined with cork and having a bifurcated body extending across it near the top of the compartment. A lure is stored in Malmquist's box by inserting the shank of its tail hook into the bifurcating slot and suspending the lure from the tail hook, the tines of which are embedded in the cork. Fabrication of the individual compartments in Malmquist's box is time consuming and expensive.\nBaumgartner, in U.S. Pat. No. 2,268,920, teaches a flexible band for mounting about the crown of a hat, the band having a plurality of safety-pin-like retainers spaced apart along its length. A fishing fly may be attached to each of Baumgartner's retainers by threading one leg of the retainer through the eyelet of the fly hook and then pushing that leg into a clamped position. Baumgartner's apparatus relies on the spacing between retainers being greater than twice the length of a lure in order to keep adjacent lures from becoming entangled. Although this may be satisfactory for relatively small dry flies, the use of Baumgartner's system for larger spinning or baitcasting lures would result either in entangled lures, or in apparatus having a severely restricted carrying capacity.\nNelson, in U.S. Pat. No. 2,364,807, teaches a lure storage box having a slotted platform. The hooks of stored lures extend downward through the slots into a separate chamber. Lures are held in position by rubber band stretched across the top of them. Nelson's apparatus involves components having complex shapes that are difficult to fabricate, and requires awkward hand motions by the fisherman who wishes to disengage a lure from the box.\nIn U.S. Pat. No. 2,065,234 Martinez teaches a fly box having a plurality of rigid rods extending across the surface of a resilient rubber sheet. A fly to be stored has its hook threaded carefully between a rod and the springy rubber sheet and is thereby held in a fixed position. Martinez's box is difficult to use without occasionally sticking the point of a hook into the rubber, which tears the rubber and diminishes its utility as a springy retainer.\nSchweigert, in U.S. Pat. No. 2,164,259 discloses storage apparatus for fishing flies, in which the curved portion of each fly's hook is draped over a rigid rod. Adjacent flies are separated by small metal tabs, and the flies are held on the rod by a single spring. Schweigert's apparatus is not suitable for the storage of multi-hooked lures, as those hooks not draped about the rod would become entangled.\nFigley, in U.S. Pat. No. 1,482,678, teaches a dry fly book in which each fly hook is retained intermediate two leaf spring elements. Figley's apparatus, like that of Schweigert and others, is not suitable for use with multi-hooked lures.\nMadsen, in U.S. Pat. No. 1,076,898, teaches a storage box for skeins of yarn and the like. Madsen provides metallic spring fingers pinching the yarn against fixed projections extending inward from a wall of the box."}
{"text": "As the display technology becomes increasingly mature, various displays have gradually developed. At present, liquid crystal displays (briefly known as LCDs) are more and more widely applied due to advantages such as little power consumption, miniaturization, light weight and thin thickness, and so on.\nAn existing liquid crystal comprises two polarizers having light transmitting axes perpendicular with each other and a liquid crystal panel located between the upper and lower polarizers. The principle of greyscale display is that the lower polarizer converts a natural light into linearly polarized light; the deflection status of liquid crystals is controlled by voltages so as to convert the linearly polarized light into elliptically polarized light, and the upper polarizer analyses polarization of the elliptically polarized light so as to realize different greyscale displays."}
{"text": "The present invention relates, generally, to the field of computing, and more particularly to call screening by applying data analytics to both digital footprint and social media.\nData analytics relates, typically, the discovery, interpretation, and communication of meaningful patterns in data areas rich with recorded information. Data analytics relies on an extensive use of computer skills, such as mathematics and statistics, combined with the use of descriptive techniques and predictive models to gain valuable knowledge from data analysis.\nA digital footprint is typically a unique set of traceable digital activities, actions, contributions and communications that are manifested on the Internet or on computing devices.\nSocial media relate to computer-mediated technologies that facilitate the creation and sharing of information, ideas, career interests and other forms of expression via virtual communities and networks of the computing devices."}
{"text": "1. Field of the Invention\nThe present invention relates to motor vehicles and in particular to a method for controlling cylinder deactivation.\n2. Description of Related Art\nMethods for controlling cylinder deactivation have been previously proposed. Bolander (U.S. patent number 2006/0130814) is directed to a method of regulating a displacement on demand (DOD) engine. The Bolander method teaches adjusting activation of a first cylinder to partially achieve the desired engine displacement and subsequently adjusting activation of a second cylinder to fully achieve the desired engine displacement. In other words, instead of activating multiple cylinders simultaneously, a first cylinder is activated, followed by a second cylinder being activated. During a first step before partial deactivation, the control device determines whether the displacement on demand system should be disabled. The displacement on demand system is disabled whenever the vehicle is in a situation where activation of the DOD system would be inappropriate. Such conditions include that the vehicle is in a transmission mode other than drive (i.e. park, reverse or low range). Other situations include the presence of engine controller faults, cold engine, improper voltage levels and improper fuel and/or oil pressure levels.\nFoster (U.S. Pat. No. 6,904,752) is directed to an engine cylinder deactivation system that improves the performance of the exhaust emission control systems. The Foster design discloses a cylinder deactivation system to control temperature and air/fuel ratio of an exhaust gas feed-stream going into an after-treatment device. Foster teaches cylinder deactivation for controlling temperature of the exhaust gas continues as long as the operating point of the engine remains below a predetermined level, or the coolant temperature is below the operating range of 82-91 degrees C., or the exhaust gas temperature is below an optimal operating temperature of the after-treatment device, e.g. 250 degrees C. In other words, the Foster device uses a single threshold limit for the engine operating level, the coolant temperature and the exhaust gas temperature.\nDonozo (U.S. Pat. No. 4,409,936) is directed to a split type internal combustion engine. In the Donozo design, the internal combustion engine comprises a first and second cylinder unit, each including at least one cylinder, a sensor means for providing a signal indicative of engine vibration and a control means for disabling the first cylinder unit when the engine load is below a predetermined value. The controller means is adapted to hold the first cylinder unit active, regardless of engine load conditions, when the engine vibration indicator signal exceeds a predetermined value indicating unstable engine operation. In the Dozono design, cylinder deactivation may occur during low load conditions any time the measured vibrations are below a particular threshold value. Dozono does not teach a method where cylinder deactivation is stopped for low load conditions based on engine speed.\nWakashiro (U.S. Pat. No. 6,943,460) is directed to a control device for a hybrid vehicle. The Wakashiro design teaches a method for determining if cylinder deactivation should be used and a separate method for determining if the engine is in a permitted cylinder deactivation operation zone. The factors used to determine if the engine is in a permitted cylinder deactivation zone are the temperature of the engine cooling water, the vehicle speed, the engine revolution rate, and the depression amount of the accelerator pedal. In each case, these factors are evaluated based on a single predetermined threshold. In other words, if each of these factors is determined to be above or below (depending on the factor) a predetermined threshold, the cylinder deactivation operation is prevented.\nWhile the prior art makes use of several parameters in order to determine if cylinder deactivation should be stopped, there are shortcomings. The prior art teaches only threshold limits above which cylinder deactivation can continue and below which cylinder deactivation should be stopped. Also, the prior art does not teach the use of stop deactivation dependent on various parameters including engine speed, vehicle speed, transmission ratio, or engine load. There is a need in the art for a system and method that addresses these problems."}
{"text": "The invention relates to a protective suit for protection against harmful chemical and biological substances as classified in the preamble of claim 1.\nPersons who, in the event of accidents involving harmful substances, work on site and therefore are exposed to these substances have to wear protective suits which stop the harmful substances from getting through to the skin or underwear.\nThe same is true of the military field when chemical or biological warfare agents are used.\nFrom commercial practice there are known for this purpose two-part protective suits consisting of a jacket and a pair or trousers, each constructed of an outer and a liner into which a carbon-based layer has been introduced. The substances adsorb onto the carbon surface. However, it is disadvantageous in this connection that the carbon will take up a multiplicity of chemicals, including chemicals which are relatively harmless, such as, for example, exhaust gases, deodorants, insect repellents and the like, and hence will become saturated, and consequently pervious, over time. Another problem here is the protective effect against chemicals in the liquid state.\nA protective effect, moreover, presupposes a certain layer thickness for the carbon, which not only creates an adverse thermal effect through a heat build-up for the wearer, but also makes the protective suit relatively cumbersome and heavy. In addition, it takes several minutes to put on the two-part protective suit and seal it, which likewise represents an appreciable disadvantage in the event of danger threatening.\nThe present invention therefore has for its object to provide a protective suit which is simpler to handle and safer in use.\nAccording to the invention, this object is achieved by the features recited in the characterizing part of claim 1.\nInstead of a bulky barrier layer of charcoal, the barrier layer which is used now is a membrane which is surrounded by a protective layer on both sides. The protective layer on the outside serves as a mechanical protector for the membrane and the inner protective layer also serves to increase the wear comfort. To this end, for example, the outer protective layer can be formed of breathable polyester. The inner protective layer can consist of polyamide which can be hydrophilicized in order that it may readily absorb perspiration and transport it to the outside. In an advantageous execution of the invention, the membrane is likewise configured to be breathable. A very useful membrane has been determined to be a cellulose-based material. Such a material is described in U.S. Pat. No. 5,743,775 for example.\nThe configuration as an coverall facilitates and appreciably speeds the putting on of the suit, since distinctly fewer locations have to be sealed. A membrane, unlike a charcoal layer, can be selectively impervious to harmful substances, while being otherwise pervious. A cellulose-based membrane blocks off a multiplicity of harmful chemicals. Moreover, no saturation occurs and at the same time protection is also provided against chemicals in the liquid state or else as an aerosol and not just protection against gas.\nThe pressure difference on both sides of the membrane ensures that moisture is transported away from the body, so that there is no heat build-up for the person who is wearing the protective suit. In addition, the evaporation processes which take place ensure a cooling effect.\nThe execution of the protective suit as an coverall provides a suit which is complete in itself, it then merely being necessary to make connecting pieces to a protective mask and to gloves.\nIn a very advantageous execution of the invention, the coverall is provided with a hood which is integrated in and hence is integral with the coverall. The same applies to the integration of socks in the coverall. This shortens the time needed to put on the protective suit even more compared with known protective suits. In addition, integration of the hood and socks in the coverall also avoids dangerous spots with regard to imperviousness.\nIn a further very advantageous execution of the invention, the at least one zipper extends beyond the hood and the hood divides into a rear part and a front part when the zipper is in the open state and hence, in the open state, can be laid down on the back and on the chest. It is merely necessary to ensure that the zipper used is impervious to the harmful substances and is also correspondingly imperviously connected to the coverall. One possible zipper of this kind is described for example in German Utility Model DE 201 10 132.7.\nAdvantageous developments and executions of the protective suit according to the invention will be apparent from the remaining subclaims and from the illustrative embodiment described hereinbelow in principle with reference to the drawing, where"}
{"text": "This invention relates to a sheet material bending angle detector and to an operation method of a press machine using the bending angle detector."}
{"text": "Battery and wall plug operated ozone generators are used to deodorize refrigerators, rooms, water, clothing, and other items."}
{"text": "Field-effect controlled transistors, such as MOSFETs (Metal Oxide Field-Effect Transistors) are widely used in automotive, industrial, or consumer electronic applications for driving loads, converting power, or the like. Such transistors, which are often referred to as power transistors, are available with different voltage blocking capabilities. The “voltage blocking capability” defines the maximum voltage level the transistor can withstand in an off-state (when switched off). The on-resistance RON of a power transistor is another relevant device parameter. The “on-resistance” is the electrical resistance of the power transistor between its load nodes (drain node and source node) in the on-state (when switched on). Typically, a power transistor includes a plurality of transistor cells that are connected in parallel. Basically, the more transistor cells the power transistor includes the lower is the on-resistance.\nThere is a need to reduce the on-resistance without increasing the size (the required chip area A) of the semiconductor body (semiconductor chip) in which active areas of the power transistor are integrated. There is therefore a need to reduce the size of the individual transistor cells."}
{"text": "Liquid processing apparatuses include a hollow-shaped bottom plate supporting a target semiconductor wafer (hereinafter referred to as a “wafer”), a rotary shaft fixedly connected to the bottom plate and rotated by a spin motor, a supply pipeline extended in the rotary shaft to supply a cleaning solution to the wafer supported on the bottom plate, and substrate lifting pins ascending to support the bottom of the wafer. One example of the liquid processing apparatuses is disclosed in Japanese Unexamined Patent Publication No. Hei 9-290197\nIn the liquid processing apparatuses, a cleaning solution, such as a chemical solution or a rinsing solution, that is used to clean wafers, may be attached to the substrate lifting pins via through holes. Accordingly, when the wafer is lifted by the substrate lifting pins and is delivered to a conveyor robot after the wafer is dried, the droplet of the cleaning solution attached to the substrate lifting pins may be attached to an opposite surface of the wafer.\nIf the cleaning solution is attached to the wafer as described above, not only watermarks may be formed on the wafer to which the droplet is attached, but also the humidity in a carrier for conveying the wafer may be increased, thereby causing the formation of watermarks on other wafers accommodated in the carrier.\nFurther, in the liquid processing apparatuses, a first path connected to an inert gas connection path and a second path connected to a cleaning solution connection path are installed inside a dual supply pipeline, and the front ends of the first path and the second path are positioned at the same height. Accordingly, the cleaning solution may flow into the first path."}
{"text": "Wireless communication systems have been widely deployed to provide various types of communication services such as voice or data. In general, a wireless communication system is a multiple access system that supports communication of multiple users by sharing available system resources (a bandwidth, transmission power, etc.) among them. For example, multiple access systems include a Code Division Multiple Access (CDMA) system, a Frequency Division Multiple Access (FDMA) system, a Time Division Multiple Access (TDMA) system, an Orthogonal Frequency Division Multiple Access (OFDMA) system, a Single Carrier Frequency Division Multiple Access (SC-FDMA) system, and a Multi-Carrier Frequency Division Multiple Access (MC-FDMA) system.\nD2D communication is a communication scheme in which a direct link is established between User Equipments (UEs) and the UEs exchange voice and data directly with each other without intervention of an evolved Node B (eNB). D2D communication may cover UE-to-UE communication and peer-to-peer communication. In addition, D2D communication may find its applications in Machine-to-Machine (M2M) communication and Machine Type Communication (MTC).\nD2D communication is under consideration as a solution to the overhead of an eNB caused by rapidly increasing data traffic. For example, since devices exchange data directly with each other without intervention of an eNB by D2D communication, compared to legacy wireless communication, the overhead of a network may be reduced. Further, it is expected that with the introduction of D2D communication will reduce the power consumption of devices participating in D2D communication, increase data transmission rates, increase the accommodation capability of a network, distribute load, and extend cell coverage."}
{"text": "Generally speaking, individuals often provide items for sale via online auction systems. Conventional online auction systems may conduct an electronic auction of the item on behalf of a provider (e.g., a seller) of the item. For example, an online auction system may make the item available for bidding by members of the general public and facilitate provision of the item from the provider and a winning bidder.\nConventional online auction systems have various limitations and drawbacks. For example, an individual supplying an item for auction may be unable to identify other individuals who may want the item and notify them that the item is available via the online auction system. Similarly, an individual who is interested in acquiring a particular item or perusing potential items of interest that are relevant to the individual may have to search a conventional online auction system and possibly wade through irrelevant items to determine whether the particular item is available or otherwise find potential items of interest.\nFurther, conventional online auction systems do not provide enough control over the individuals who may participate in the auction. For example, while some online auction systems may allow the individual providing the item to block particular individuals from participating in the auction, they typically do not facilitate targeting potentially interested individuals to participate in the auction.\nConventional online auction systems suffer from these and other drawbacks."}
{"text": "The present invention relates generally to a memory array address decoder and more particularly to a memory array address decoder used by a dynamic random access memory device.\nA typical DRAM memory device is comprised of a plurality of memory cells, each comprised of a transistor and a capacitor. Each memory cell stores one bit of data in the form of a voltage. A high voltage level (e.g., 3V) represents a logic “1”, whereas a low voltage level (e.g., 0V) represents a logic “0”. The DRAM may also include peripheral devices, such as drivers, sense amps, input/output devices, and power supplies, etc., that are used to identify memory cells, access the memory cells, and store information within and read information from the memory cells, among others.\nThe memory cells may be arranged in an array with each memory cell being connected to a wordline and a digitline. Each memory cell has a unique address. Typically, the DRAM's control logic receives commands (e.g., read, write, etc.) and address information from a memory system controller. The address information is decoded by row and column decoders and the specific memory cell for which the command is directed is identified and the command executed.\nMost DRAMs have built-in redundancy. Thus, should a memory cell become inoperable, a redundant memory cell can be assigned to logically take its place. For DRAMs incorporating redundancy, address information is sent to a comparator circuit. If the address information corresponds to an inoperable memory cell, a match signal is generated which substitutes the memory address of a redundant cell for that of the inoperable memory cell. The generation of the match signal, however, limits the speed of the DRAM because the other circuitry within the device must wait for the match signal to be generated.\nThe address information supplied to the row and column decoders can be either “true and complement” or pre-decoded as is known in the art. Pre-decoded address lines, for example, may be formed by logically combining (i.e., using one or more AND logic gates) true and complement addresses. Pre-decoded addressing requires less power than true and complement addressing because fewer signals need to make transitions during address changes. Additionally, pre-decoded addressing has a higher efficiency than true and complement addressing because fewer transitions are required to decode the same number of addresses.\nTypical prior art decoders are classified as static or dynamic. FIG. 8 is a schematic of a static column decoder 80 according to the prior art. The static column decoder 80 includes NOR gate 81, NAND gates 82 and 84, and inverters 83 and 85. Pre-decoded address signals CA345i<0> and CA67Ei<0> are input into NOR gate 81 (where the “i” indicates that the signal is active low). The output of NOR gate 81 is fed to an input of NAND gate 82 and to an input of NAND gate 84. The pre-decoded address signal CA012<1> is provided to another input of NAND gate 82, whereas the pre-decoded address signal CA012<0> is provided to another input of NAND gate 84. The output of NAND gate 82 is fed to the input of inverter 83, which outputs the column select signal CSEL<1>. The output of NAND gate 84 is fed to the input of inverter 85, which outputs the column select signal CSEL<0>.\nThe static column decoder 80, although simple to implement, has several deficiencies. First, different gate delays are created for each column select line because the static column decoder 80 uses CMOS gates. Although originally tolerable, the differing gate delays create problems for today's DRAMs which operate at increased speeds. Second, the static column decoder 80 requires that the turn-on period and turn-off period of a match signal, for example generated when the originally addressed cell is found to be inoperable, be equal (i.e., the rise and fall times of the match signal must be the same). Thus, the cycle time for a static column decoder 80 is adversely affected.\nFIG. 9 is a schematic of a dynamic column decoder 90 according to the prior art. The dynamic column decoder 90 includes several p-mos transistors (M2, M4, M6), several n-mos transistors (M1, M3, M5, M7), and several inverters (91–96). Typically, the dynamic column decoder is operated by first applying a precharge signal to the gates of transistors M3 and M7, thus pulling nodes 98 and 99, respectively, to ground. The precharge signal is then removed from the gates of transistors M3 and M7. Inverter 91 is used to latch node 98 at ground (thus, preventing node 98 from floating when the precharge signal is removed from the gate of transistor M3). Likewise, inverter 94 is used to latch node 99 at ground (thus, preventing node 99 from floating when the precharge signal is removed from the gate of transistor M7).\nAfter the precharge signal is removed, the pre-decoded address signals CA012i<1>, CA012i<0>, CA345i<0>, and CA67E<0> are applied to the dynamic column decoder 90. The state of each of the output signals CSEL<1> and CSEL<0> is dependent upon these pre-decoded signals as should be apparent to one skilled in the art.\nOne advantage of the dynamic column decoder 90 over the static column decoder 80 is that the dynamic column decoder has consistent gate delays. Thus, the column select output lines (i.e., signals CSEL<1> and CSEL<0>) have consistent on/off times. The dynamic column decoder 90, however, has several deficiencies. For example, the feedback inverters 91, 94 must be sized large enough to keep the nodes 98 and 99, respectively, from floating (i.e., must keep the nodes 98 and 99 at ground potential, set when the precharge signal goes active), but sized small enough to be easily overridden when the pre-decoded address signals CA012i, CA345i, and CA67E go active. Additionally, the pre-decoded address signals CA012i, CA345i, and CA67E cannot go active until the precharge signal goes inactive, else the pre-decoded address signals overlap and fight the precharge signal. If the precharge and pre-decoded address signals CA012i, CA345i, and CA67E overlap, the turn on timing of the dynamic column decoder is adversely affected. Thus, the precharge signal consumes a relatively large amount of the cycle time which is available for providing the column select signals.\nThus, there exists a need for an improved memory array decoder that has a consistent turn-on time, that can utilize a precharge signal without adversely affecting the cycle time, and which overcomes the other limitations inherent in prior art."}
{"text": "1. Technical Field\nAn aspect of the present invention relates to a computer readable storage medium for a computer capable of displaying operational status information concerning a device such as a printer.\n2. Related Art\nA device driver, which is a computer program to be run on a personal computer (PC) to monitor and manipulate behaviors of a device (e.g., a printer) connected to the PC, is known. The device driver receives operational status and conditions of the device and displays information concerning the status through a display of the PC."}
{"text": "The present invention relates to a drive train cooling arrangement for motor vehicles, having a first cooling circuit and a second cooling circuit and having a pump arrangement, by means of which coolant can be fed to the first and the second cooling circuit.\nThe present invention furthermore relates to a method for operating a cooling arrangement of this kind. There is a known practice in drive trains of motor vehicles of feeding cooling or lubricating oil to subassemblies of a transmission, such as bearings and gearing, by way of injection lubrication. In this case, a pump is generally provided in the region of a transmission sump, and this pump can be driven by means of a drive motor such as an internal combustion engine, for example.\nIt is furthermore a known practice in the area of drive trains to cool wet multiplate clutches of the kind often used in dual clutch transmissions, for example. In this case, a separate pump is used and is generally likewise coupled to a crankshaft, building up a pump pressure continuously in operation.\nTo cool such friction clutches, use is generally made of an automatic transmission fluid of relatively low viscosity, whereas “hypoid” oil with a higher viscosity is used for lubrication in many transmissions.\nHere, a volume flow of such a cooling or lubricating fluid is generally supplied via electromagnetically controlled valve arrangements. Valves of this kind generally contain spool valves, which make great demands as regards cleanliness during manufacture and assembly.\nIn the case of drive trains for motor vehicles which additionally have an electric motor as a drive motor (hybrid drive trains), it may also be necessary to cool the electric motor. For this purpose, another cooling arrangement is provided, and this is likewise generally operated with automatic transmission fluid. In this case, an electrically driven coolant pump is generally provided to supply the electric drive motor with coolant in order to allow a purely electric drive mode as well."}
{"text": "State-of-the-art magnetic field-based current sensors generally include a magnetic field sensor, an iron core, and an amplifier. The iron core surrounds an electrical conductor, preferably comprising copper for carrying the current to be measured, the current causing a magnetic field corresponding to the intensity of the current to be generated in the iron core. The magnetic field sensor is arranged in a gap in the iron core. Thus, the field lines of the magnetic field induced in the iron core by the current to be measured pass through the magnetic field sensor. The magnetic field sensor generates an electrical output signal generally proportional to be current to be measured, and this signal is amplified by an amplifier, such as an operational amplifier.\nThe iron core in conventional current sensors primarily provides two functions. First, the conducting of the magnetic field induced by the current to be measured through the magnetic field sensor, and second, the screening of the magnetic field sensor from extraneous electric and magnetic fields which are known to cause measurement error. However, the use of iron cores results in a sensor cost that is too expensive for many applications.\nA known improved current sensor design disclosed in EP 1192473/WO 00/79291 dispenses with the iron core and thus its expense. According to the sensor design disclosed, the primary conductor which carries the current to be measured comprises an electrically conducting ferritic material generally mixed with an electrically conducting material (e.g. copper). The primary conductor at least partially surrounds the magnetic field sensor which assumes the functions of the iron core. The magnetic field sensor is centered near a notch in the primary conductor which serves to focus the magnetic flux lines from the primary current in the primary conductor that reaches the magnetic field sensor. The current sensor design disclosed in EP 1192473/WO 00/79291 generally has the disadvantage that the electromagnetic field associated with the current to be measured generates a magnetic field in the primary conductor, which can prevent precise measurement of the magnetic field associated with the current to be measured in the notch by adding another magnetic field component. The remanance of the primary conductor also causes non-linearity and magnetic saturation in the region of the notch which degrades the effect of the focussing notch. If only magnetic permeable material is used for the primary conductor (without copper or other electrical conductor), then a large cross-sectional area of this expensive material is generally needed because of its relatively high electrical resistivity.\nThus, although sensors disclosed in EP 1192473/WO 00/79291 having a primary conductor including magnetic permeable material mixed with electrically conducting material can provide good performance and a reduced cost compared to conventional iron core-based current sensors, in certain important applications more precise measurement of electric current is required. Precise measurement of electric current requires improved linearity (between the current to be measured flowing as primary current and the magnetic flux reaching the magnetic sensor) and reduced resistive impedance. Lower cost would also be helpful.\nSome present current sensor applications include measurement of electric current include motor current control in frequency inverters, phase current control in servo-drives, overload protection for motors and power semiconductors, control of energy systems (e.g. wind energy), current control and monitoring in welding equipment, current regulation in switching mode power supplies, current control in uninterruptible power supplies, battery current diagnosis in motor vehicles, ground fault detection and laboratory test equipment. Some of these applications generally require precise current measurements.\nSmall sensor size is required in certain applications, such as those where current is measured directly on a printed circuit board. Moreover, low cost is also required in most automotive applications and other high volume applications. What is needed is a new current sensor design that provides improved sensor performance while also providing a low cost."}
{"text": "1. Field\nThe present disclosure relates to an adhesive tape with print that is affixed to an adherent and used, and an adhesive tape roll with print that uses the same.\n2. Description of the Related Art\nThere are known adhesive tapes with print that are affixed to an adherent and used. These adhesive tapes include a tape base layer (resin film base), an adhesive layer, and a separation material (separation sheet), in that order, from the front surface side to the rear surface side in the thickness direction, and print formation is performed by transferring the ink of the ink ribbon to the front surface side. The ink ribbon used is a ribbon comprising a ribbon base layer and an ink layer that sticks to the front surface side of the adhesive tape. The user uses the adhesive tape of a preferred length by affixing it to a suitable adherent.\nWhen the user affixes and uses the adhesive tape as described above, dirt and debris sometimes stick to the front surface, and it is contemplated that a separation layer could be provided on the front surface of the tape base layer to prevent this. Hence, such an adhesive tape may be conceivably generated as a roll wound around a predetermined axis, for example. In a case where a separation layer is thus provided, the adhesive tape includes a separation layer, a tape base layer, and an adhesive layer, in that order, and thus the adhesive layer in the roll is in contact with and adhered to the separation layer in the radial direction of the roll. Then, when the adhesive tape is fed out from the roll, the adhesive layer is sequentially peeled and separated from the separation layer, requiring consideration for improving the separability thereof as well. Furthermore, in a case of a structure that provides separability as described above, consideration must be made to improve the adhesion between the separation layer and ink ribbon in order to ensure favorable transferability at the time of transfer from the ink ribbon.\nIn the prior art described above, achieving an adhesive tape with optimum characteristics that is capable of maintaining separability at the time of feed-out from the roll, transferability from the ink ribbon, and antifouling properties when affixed and used as previously described was not taken into account."}
{"text": "Shoppers in retail stores typically use shopping carts or baskets owned by the store to hold their purchases as they travel through the store. Purchases are typically bagged at the check out counter and either reloaded into the cart or the bags themselves are carried out of the store by the shopper. Carrying the bags is cumbersome for the shopper and the number of bags that a shopper can carry is limited by the strength of his or her arms. With regard to using shopping carts outside of the store, such an arrangement does not work for city shoppers who walk to and from the store. In addition, even if the shopper drives to the store, and carries the groceries to his or her car with a store-owned cart, the store must go to the trouble of collecting the carts from the parking lot. In addition, the carts must be maintained, repaired and stored, which can add up to significant costs and take up valuable storage space.\nConsumer-owned folding wire carts for use in shopping are known. The shopper brings the cart to the store, opens/unfolds it, and uses it to hold items during shopping. The shopper then uses the same cart after checkout to transport the purchases to his or her car or home. The disadvantage of such carts, however, is that they can be heavy due to the metal wire construction. In addition, they tend to feature a very utilitarian appearance. They also typically do not fold completely flat, and thus take up valuable space which is undesirable (especially when the shopper is an apartment dweller).\nIn view of the above, a need exists for a collapsible personal trolley that is inexpensive, lightweight, durable and attractive."}
{"text": "Although data indicate that the immune system is of major importance for cancer control, malignant tumours continue to grow and the efficacy of immunotherapy is rather poor with an objective remission rate of 15-30%. There can be several reasons for this apparent paradox:                Tumours avoid the recognition by the immune system by not expressing tumour associated antigens properly        Tumour associated antigens (often self antigens), which are too weak to elicit an adequate immune response        Induction of tolerance        Cancer related immunosuppression, which prevents an adequate immune response        \nThese alternatives require completely different therapeutic strategies, either proper stimulation of the immune system or control of cancer related immunosuppressor mechanisms.\nImmunosuppression in cancer is mainly characterized by: Reduced proliferative and cytotoxic capacity of lymphocytes, in particular tumour infiltrating lymphocytes, poor migration of inflammatory cells, reduced production and response to IL-2, difficulties to elicit an immune response by vaccination, also against other than tumour related antigens and pathological cytokine production. This dysregulation of the immune system results in poor immune mediated cancer control and a paraneoplastic syndrome (subfebrility, fatigue, anorexia, weight loss and deterioration of laboratory parameters).\nImmunostimulatory therapeutic strategies using cytokines (e.g. interferons, interleukins) or vaccination, in order to enhance the immune mediated reactivity to the tumour, has been tried for several decades, but have so far had only very limited success. This indicates that immunostimulation in order to overcome a poor immune response in cancer patients might be suppressed by other so far unidentified mechanisms.\nWe have in two previous patent applications described two fundamental immunoregulatory mechanisms of relevance to all types of malignant tumours. In the first of these applications the importance of Fc-receptor modulation and ways to overcome this cancer related immunosuppression by modulating Fc-receptor cross-linking was demonstrated. In this patent application also proteolytic fragments of normally occurring proteins were demonstrated to induce pathological monokine production. In the second patent application some of these neo-structures were found to be integrin binding/blocking and their occurrence and immunoregulatory activity was further analysed by using monoclonal antibodies directed against albumin derived neo-structures. In the latter patent application also the occurrence and importance of auto-antibodies to these neo-structures is described."}
{"text": "It is already possible using artillery locating radar and other surveillance systems, for example, to determine rapidly and with high precision the location of an artillery gun that has opened fire. There is thus a good opportunity for an enemy to open effective counter-battery fire. The artillery has therefore more or less been forced to depart from its previously fairly stationary tactics in favour of significantly more mobile tactics involving rapid engagements in the form of short intensive fires followed by immediate redeployment to a pre-determined deployment site at a sufficiently safe distance from the previous one. These new tactics have resulted in an increased need for every gun to be self-propelled and capable of carrying at least a primary requirement of ammunition.\nOne must also assume that coming generations of artillery will use modular type propellant charges, that is, propellant charges consisting of a number of modular charges of different sizes, such as length and, to a certain extent, diameter, of different charge strength with primarily rigid combustible outer casings, and that are combinable in various ways to provide the desired muzzle velocities. At present this system of modular charges is called M(A)CS, that is, Modular (Artillery) Charge System. Moreover, as the next generation of artillery guns is expected to be equipped with armoured protection against battlefield fragments to an even greater extent than is normal today. Next generation loading systems will be required to operate very rapidly and be capable of stowing large quantities of propellant charges and of handling all the different types of propellant charges in the M(A)CS. The propellant charges must also be stowable in the least possible space. In addition, loading systems shall be robust and durable, and the propellant charge magazine shall be replenishable in a very short time, preferably from a vehicle equipped with an automatic resupply unit."}
{"text": "Today metal walls are used in a wide variety of buildings including both the residential and commercial buildings. Basically, a metal wall section includes a series of spaced apart metal studs that extend between an upper and lower tracks that ordinarily are secured to an underlying floor and an overhead support structure. In some cases, the upper track may be free floating. Typically secured to the sides of the studs is a wall covering such as gypsum board. Most metal walls are straight. However, in recent years it is noted that many architects and interior designers have started to use curved walls in particular locations in buildings. In some cases these curved walls may be used purely for aesthetic or design appeal. In other cases, they may serve utilitarian functions.\nIn straight metal walls, the upper and lower tracks typically comprise an elongated U-shaped metal channel. Such channels are, of course, not suited for curved wall constructions. Yet, in order to fabricate a curve wall, it is necessary to provide a support or a retaining structure that runs along the upper and lower ends of the studs that comprise the wall.\nIn the past, there have been attempts at providing flexible track designs. For example, reference is made to U.S. Pat. No. 6,000,181 and the prior art discussed in this patent. U.S. Pat. No. 6,000,181 is expressly incorporated herein by reference. This patent discloses a multi-section track wherein the respective tracks are pivotally connected with respect to each other. However, this flexible or deformable track design includes an adjustable strap that is threaded through the upstanding stud tabs formed on each section. These straps make the flexible track complicated and expensive to manufacture. This complexity also can have an adverse affect on the ease of use of this flexible track.\nTherefore, there has been and continues to be a need for a flexible or deformable track for a curved metal wall that is relatively simple in construction and which is easy to use."}
{"text": "The present invention relates to a video display device such as a TV receiver or a image processing apparatus such as a set-top box including a DVD player or a digital tuner, and more particularly to an image processing apparatus with a configuration in which the luminance distribution of an input video signal is detected for correcting the image quality.\nA prior art technology for improving the contrast effect of a video displayed on a display device is disclosed, for example, in JP-A-2002-359754. The disclosed technology detects a luminance histogram from the input video signal and corrects the minimum value of the video signal according to the distribution amount of the minimum interval range in the luminance histogram for adjusting the gray scale near the black."}
{"text": "1. Field of the Invention\nThe invention generally relates to network traffic analysis and in particular, to click fraud detection.\n2. Description of the Related Art\nAdvertising is one of many ways to direct Internet traffic to a website. Advertisers or advertising sponsors typically purchase online advertising, such as sponsored links or sponsored ads, and pay for such advertising based on, for example, a block of advertising impressions, a per click basis, conversions, leads, actions, or the like. These online advertising services can be subjected to fraud.\nPay per click (PPC) advertising is an arrangement in which a publisher displays clickable links from an advertising sponsor in exchange for a charge per click. An advertising network acts as a middleman between these publishers and advertising sponsors. Typically, for each click on an advertisement (ad), the advertising sponsor pays an advertising fee to the advertising network, which in turn pays the publisher a share of this money. This revenue sharing system can be exploited for click fraud.\nRelatively large advertising networks, such as Google's AdWords/AdSense and Yahoo! Search Marketing, can also be considered to be publishers (on their search engines and their various websites). This complex relationship can create a conflict of interest. For example, an advertising network would pay advertising fees to a publisher for undetected click fraud, but the advertising network would also collect advertising fees from an advertising sponsor. Typically, the advertising network charges more to the advertising sponsor than it pays out to publishers, so there is typically little incentive for the advertising network to detect fraudulent clicks.\nAdvertising networks can attempt to stop fraud, but determining which clicks are legitimate can be a difficult task. Advertising sponsors typically resist paying for fraudulent clicks. Publishers typically resent having to pay refunds for click fraud detected late.\nClick fraud occurs in pay per click (PPC) online advertising when a person, automated script, computer program, or the like, imitates the clicks of a legitimate user of a web browser to generate an improper charge per click.\nThose engaged in large scale fraud will often run scripts that simulate a human's clicks on ads in web pages. However, huge numbers of clicks appearing to come from just one, or a small number of computers, or a single geographic area, can look suspicious to an advertising network and to advertising sponsors.\nSophisticated scammers circumvent detection by clandestinely controlling relatively many computers with their own Internet connections running in disparate geographic locations. Often, scripts do not truly mimic human behavior, so scammers can turn to Trojan code or other malware to turn an unsuspecting person's machine into a zombie computer and use sporadic redirects or DNS-cache-poisoning to turn the user's actions into clicks on PPC advertisements. Such actions can generate revenue for the scammer, who may be paid by a beneficiary of click fraud.\nAn example of a conventional solution typically only detect click fraud in pre-defined advertising campaigns, and advertising sponsors are alerted to multiple page loads, or clickthroughs from identical Internet Protocol (IP) addresses. Present solutions can typically examine activity only on a URL by URL basis for the HTTP referrer resource, which can be cumbersome when many URLs are providing referrals to the referred-to website."}
{"text": "This invention relates to the field of computer systems. More particularly, a system and methods are provided for inlining across protection domain boundaries within a system virtual machine.\nA system virtual machine (VM) executes instructions from a source or “virtual” instruction set architecture (vISA), such as the x86 ISA, on a “native” ISA (nISA), such as SPARC®, by interpretation or dynamic compilation. Use of a virtual machine thus allows one physical computing device to appear to host one or more virtual computers, which can execute a variety of operating systems and application software.\nA VM awards a system designer great freedom in design within the nISA. In particular, in traditional non-virtual machine computer systems a processor executes instructions essentially one at a time, using only local knowledge of the current process. A VM, however, can apply optimizations across greater regions of a program, including inlining, whereby instructions from the target of a function call are compiled into, and optimized with, the source of the call. Currently, however, inlining can only be performed for function calls fully contained within the boundaries of one protection domain—a combination of privilege level and applicable address space.\nHowever, many application workloads involve a wealth of system calls (e.g., calls from application code into the operating system), especially for input/output (I/O). For example, in some hypervisor-based systems (e.g., Xen®), I/O follows a complicated path from an application to the operating system's stub device driver via a system call, then via a hypercall to the real device driver running in another guest OS instance.\nAll these protection domain transitions (e.g., user to system to hypervisor to system) are expensive on conventional (non-system virtual machine) hardware. For example, in response to a trap or software interrupt from application code, among other operations some or all of the following actions occur: a processor's pipeline is flushed, control registers are modified, internal state (i.e., context) of the processor is reconfigured to raise the privilege level to supervisor mode from user mode, information about the user context is saved, and control is transferred to an appropriate service routine. Yet, abstractly, system calls/hypercalls are similar to function calls except that they call across a protection domain boundary.\nSimilarly, inlining cannot be performed when a call is made from a process executing in one ISA to a process executing in a different ISA. For example, a given VM may host guest operating systems that differ from each other in terms of their ISA (e.g. SPARC, x86), but which must cooperate as described above to share a real device driver. Not only do the protection domain boundaries present obstacles, but the difference in vISAs further frustrates efforts to inline calls between the guest operating systems."}
{"text": "1. Field of the Invention\nThe present invention, in general, relates to gimbals and, more particularly, to permanent magnet wide-gap motors that are used to move the gimbal payload about its inner axes.\nGimbals are commonly used to hold sensors stable when mounted on a moving vehicle, be it a land based vehicle, a sea (i.e., a water based) vehicle such as a boat or ship, or an air based vehicle such as an airplane.\nThe ability to hold a sensor stable while the vehicle moves is useful for a great variety of purposes. These purposes include obtaining information useful for navigation. Another purpose relates in general to an ability to align and then to hold the sensors where desired. All manner of sensors may be used, for example, television cameras as well as other sensors that use any preferred technology. Whatever information is being provided by the sensors is more reliable if the sensors themselves are held steady.\nIn general, gimbals have a plurality of outer axes and a plurality of inner axes. Course adjustments are commonly accomplished by movements made along the outer axes. Finer adjustments are commonly made with the inner axes.\nThere are a number of discreet functions a gimbal must achieve. It must both properly orient, maintain position, and support the size and weight of the sensors. This can vary from application to application.\nThe sensors are placed inside of a gimbal shell along with numerous other component parts that are used to orient the gimbal ball as required. The gimbal shell is often in the general shape of a ball and is therefore sometimes also referred to as a gimbal ball or simply a ball. In general, for any given size of the gimbal shell, the space (i.e., volume) that is available for the sensors is limited and a more usable volume is desirable. Accordingly, it is desirable to locate the motors that move the payload within the ball (i.e., the inner axes) as far away from the center payload area as possible. Another way of stating this objective is to locate the motors in a space that is not well used by sensors.\nGimbals often include a cardan assembly that is disposed within a ball. The cardan assembly supports the weight of the payload that is carried by the gimbal as well as allowing small rotational movements (inner axis motion) that affect the positioning of the payload within the ball.\nThese changes in position are accomplished by rotating the payload (within the gimbal ball) about three axes (typical), namely elevation, roll, and azimuth. Courser adjustments are accomplished by moving the gimbal ball itself typically about the two outer axes, elevation and azimuth.\nThe cardan assembly includes a cardan shaft that spans the inside diameter of the ball. The center of the cardan shaft is used, in certain designs, to define the internal elevation axis.\nPrior art designs place the cardan shaft so that it also aligns with the external axis. In particular, the internal elevation axis is set to align with the external elevation axis.\nThe instant invention is applicable for use with prior art designs and also with a newer offset cardan type of gimbal, as is described in an application filed for an Offset Cardan in the section, “Related Applications”.\nThe sensor(s) is used to look at any given object(s) of interest, which is sometimes called a target(s). For example, the gimbal may be disposed in an airplane and it may be used for mapping terrain or for some other purpose. There may be a reference location that the gimbal must periodically slew to, perhaps to periodically coincide the area that is being mapped with a reference location or perhaps to periodically look for changes that might be occurring at the reference location. Then the gimbal must slew back to its original orientation to continue the mapping process. Ideally, this slewing action is preferably accomplished as quickly as possible.\nThis type of slewing typically occurs primarily along the azimuth axis. Smaller angular changes typically occur on the pitch and roll axes. Another common name applied to motion about the azimuth axis is “yaw”. This is when an airplane, for example, moves left or right as controlled by its rudder. Such motion in an airplane is referred to as yaw and it occurs about the azimuth axis.\nBecause the payload has mass, the ability to slew quickly, especially about the azimuth axis, is related to the torque that is applied. A greater torque for any given configuration results in faster slew times, which is desirable.\nThere is another problem that limits the slew rate along the inner azimuth axis for a gimbal. The outer azimuth axis serves to keep a viewing window disposed in front of the sensor (i.e., camera) within its range of travel, which includes a greater range of travel than the inner axis. However, the inner axis must also slew as fast as the outer axis in order to prevent the limited rotation inner axis from exceeding its travel, that is to keep the sensor (camera) pointing at its intended target.\nIf the inner axis cannot accelerate as fast as the outer axis, the rate of slew for the inner axis will slow down and therefore limit the slew rate for the outer axis. This makes it especially important to be able to accelerate the inner axes as fast as possible.\nTherefore it is desirable to speed up the inner rate of slew, primarily along the azimuth axis because that axis is where the greatest angular degree of motion and also where periodic changes in orientation often occur.\nAnother problem relates to the shaking of the payload during use of the gimbal. Isolators (i.e., any means to isolate the payload from the gimbal ball and the vehicle upon which the gimbal is itself mounted) are provided in virtually all high performance gimbals to isolate vibration and other changes in position of the vehicle, for example the airplane, apart from that of the payload (i.e., the sensors). The object is to minimize the sensor's vibration as it is pointed at the intended object of interest while the airplane vibrates, for example.\nThe isolators allow the payload to translate (i.e., move) within a given range of motion inside of the ball. Wide-gap motors provide a relatively wide space intermediate the magnets thereof with uniform magnetic flux. This allows space for the payload to move without impacting the magnet or any other part of the motor as the isolators translate.\nIt is important to understand that the internal axes provide finer corrections than do the external axes and accordingly, a smaller range of motion is therefore acceptable for the payload in the gimbal ball. Larger corrections are made by moving the entire gimbal ball relative to the vehicle upon which the gimbal itself is mounted.\nIt is useful to note that the cardan assembly may be used to support multiple types of sensors simultaneously as the payload. For example, a zoom television camera can be used for general spotting purposes and to locate an object of interest as well as for general pointing (i.e., aiming) of the gimbal. Upon locating the object of interest, a larger focal length camera can be used to more carefully study it. Accordingly, both types of cameras can be simultaneously mounted as part of the payload that is supported by the cardan assembly.\nThe payload may also be active instead of passive. A passive payload merely observes the object of interest whereas an active payload is adapted to affect it. The payload may be used to support an active component that can, for example, illuminate the object. For example, a gimbal can contain a source of illumination, such as a spotlight or a laser, and be mounted on, for example, a helicopter. Accordingly, as the helicopter hovers and fluctuates in its position relative to the object, the gimbal can be used to compensate for any movement by the helicopter in order to hold the source of illumination constantly upon the object.\nIf the source of illumination is a spotlight, then a larger physical payload capacity allows for a larger and brighter spotlight to be used. A larger payload typically increases the mass that is being carried which tends to slow down the rate of slewing. It is desirable to optimally locate as high a torque motor as possible which, all other factors being equal, allows for a greater payload mass and therefore increased utility for the gimbal.\nAccordingly, there exists today a need for a cross plane wide-gap motor system for a gimbal that improves the aforementioned prior art limitations. In particular there is a need to step with high torque and allow good stability with the same motor arrangement simultaneously.\nClearly, such an apparatus would be useful and desirable.\n2. Description of Prior Art\nGimbals are, in general, known. While the structural arrangements of the known types of devices, at first appearance, may have similarities with the present invention, they differ in material respects. These differences, which will be described in more detail hereinafter, are essential for the effective use of the invention and which admit of the advantages that are not available with the prior devices."}
{"text": "There has been a great deal of attention directed towards the elaboration of organic compounds from readily available building blocks. One reason for such attention relates to the synthesis of pharmaceutical molecules, synthetic polymers, pesticides, food supplements, and a host of other, useful materials. Moreover, the preparation of electronics chemicals can benefit greatly from improved synthetic schemes. Of particular importance is the coupling of compounds between their respective carbon atoms, which is an important tool for synthesizing and building organic compounds. The coupling of compounds is usually complicated due to sensitivity of different functional groups on the compounds to the reaction conditions or catalysts. While some methods for coupling compounds are available, they do not always offer the efficiencies or economies required for large scale production and are not appropriate for all syntheses. Accordingly, there is continuing interest in improving methods of synthesis involving carbon-carbon coupling, called cross-coupling when between different organic molecules.\nThe art of cross-coupling reactions catalyzed by palladium catalysts is a powerful method to combine organic molecules by joining their carbon atoms. In the Suzuki cross-coupling reaction, carbon atoms of a boronic acid and an aryl halide are coupled via palladium catalysis under basic conditions. Recently, palladium catalysts, such as Pd(PPh3)4, have been used with the Suzuki reaction to cross-couple an aryl or a vinyl boronic acid to an aryl or vinyl halide or triflate. Certain other palladium catalysts, like PdCl2(dppf) and Pd(OAc)2, are also available. Primary alkyl groups can also be transferred by the Suzuki reaction using 9-BBN reagents and PDCl2(dppf) as a catalyst.\nA category of palladium catalysts useful in cross-coupling reactions are palladium-phosphine complexes. In one example, Sunagawa et al., U.S. Pat. No. 5,216,186, describe the synthesis of crystalline palladium tetrakis(triphenylphosphine) or Pd(PPh3)4. Sunagawa et al. also describe palladium tetrakis(triphenylphosphine) to be useful in the field of organic synthesis, explaining its use as a catalyst for deprotecting reactions of allyl esters or allyl ethers or allylation reactions of carbon or nitrogen atoms. Beller et al. in U.S. Pat. No. 5,831,107 describe a number of palladacycles and their importance as catalysts for processes such as the synthesis of substituted styrenes, preparation of stilbenes and cinnamic acids from aryl halides.\nDespite the usefulness of palladium-phosphine complexes and their use in cross-coupling reactions, the use of phosphine ligands often can result in side reactions such as phosphonium salt formation and aryl-aryl exchange between substrate and phosphine. In addition, it is known that highly sensitive compounds often do not tolerate the basic conditions required by the Suzuki reaction. Also, this reaction requires extensive refluxing of the reaction mixtures. Furthermore, a popular phosphine catalyst, tri-tert-butylphosphine is known to be very expensive and may not be practicable for large scale production of materials.\nAccordingly, there still remains a need for an economical, selective and efficient means for cross-coupling of compounds, particularly aryl compounds, via their respective carbon atoms. Further, there is a need for an economical, selective and efficient cross-coupling reaction that provides mild conditions for the synthesis of aryl derivates such as arylpyridines and aryl anilines, particularly sterically hindered ortho-substituted anilines."}
{"text": "Today, people often utilize computing devices for a wide variety of purposes. Users can use their computing devices, for example, to communicate with other users. Such communications are increasingly popular over a social networking system. Digital communications, such as those on a social networking system, may involve various types of communication. Some types of digital communication allow a user to engage in focused exchanges. For example, the user may target a particular user or users through the use of a messaging system or an email system supported by a social networking system. As another example, the user can enter into audio communications or video communications with other users.\nIn many instances, video communications can be preferred by users because video communications can allow the users to most effectively convey information and simulate real life communications. In some instances, two participants in different locations can engage in video communications. It also can be desirable to allow a group of users in multiple locations to use video communications to facilitate communications among the group."}
{"text": "This invention relates generally to pneumomassage apparatus adapted to displace excess fluid from the body part of a patient and more particularly, to an articulated sleeve for this purpose which fits onto and conforms to the foot and lower leg of the patient.\nLymph is a clear fluid which circulates in tissue spaces of the vertebrates and by way of a tubular network passes into the venous system. Lymph is derived from the liquid plasma of the blood but without the red corpuscles. Lymphedema is a disorder that often follows a surgical procedure such as a lymphnode dissection of the groin in the treatment of cancer. It results in an excessive accumulation of lymph in the body tissues which if not reduced may have serious consequences.\nTo reduce lymphedema or other abnormality giving rise to excess fluid in an afflicted body part, such as an upper or lower limb, it is known to provide an apparatus to compress successively portions of the afflicted part to produce a sequential or peristaltic action pumping excess fluid toward the heart. This apparatus can also be used to improve blood circulation.\nU.S. Pat. No. 4,338,923 discloses an inflatable cell-type apparatus for treating edema, the apparatus including a sleeve of flexible material divided into a plurality of internal inflatable cells extending along one dimension of the sleeve, each cell having a port for insetting and outletting air to individually inflate or deflate the cells. The sleeve when applied to the body part to be treated then surrounds it with inflatable cells extending annularly around the sleeve. Means are provided to apply pressurized air to the ports in accordance with a predetermined sequence for inflating and then deflating the cells.\nIn the pneumomassage apparatus disclosed in U.S. Pat. No. 5,014,681 for applying intermittent compression to a body part, use is made of an inflatable sleeve divided into successively overlapping inflatable cells. Pressurized fluid is applied cyclically to successive groups of cells, so as successively to inflate each group while at the same time partially deflating the preceding group.\nA useful feature of the sleeve shown in U.S. Pat. No. 5,014,681 is that the sleeve is formed from a trapezoidal blank, one inclined edge of which is provided with one component of a zipper. Adjacent the opposing inclined edge is a row of spaced complementary zipper components. To fit the sleeve to the limb of a patient, the zipper component on the one edge is connected to a selected complementary component on the opposing edge, thereby creating a sleeve whose dimensions conform to the limb of the patient.\nAlso of prior art interest is the pneumomassage device disclosed in the U.S. Pat. No. 4,374,518 in which the device is provided with a boot having a series of compartments which conform to the contour of a human foot and leg, the foot part of the boot. A compressor is provided for successively inflating and deflating the compartments of the boot in a rhythmic, preselected cycle.\nThe inflatable cell-type body treating apparatus disclosed in U.S. Pat. No. 4,338,923 for the treatment of edema includes a flat inflatable band divided into a plurality of internal inflatable cells extending along one dimension of the band, and wrappable about the body part to be treated to form a sleeve with the inflatable cells extending annularly around the sleeve. The inflatable cells are in partially overlapping relationship widthwise of the band.\nInasmuch as a sleeve in accordance with the invention includes overlapping cells inflated and deflated sequentially by a controlled compressor in a manner similar to those of the pneumomassage sleeves disclosed in the above-identified U.S. prior art patents, the disclosures of these patents are incorporated herein by reference.\nIn order to massage both the foot and lower leg of a patient afflicted with edema, it is known to fit a conical or tubular sleeve on the leg and to fit a separate sleeve on the foot. The advantage of this arrangement is that the patient, as he is being treated, is free to maintain his foot in a natural position and to flex it. This freedom of foot movement promotes the comfort of the patient and is desirable when the treatment is to be carried out for a prolonged period. However, the need for separate foot and leg sleeves and to sequentially inflate the cells in the sleeves introduces unwanted complications.\nOn the other hand, where the apparatus for treating a patient is in a boot format as in the \"\"518 patent, it can then only be used with those patients whose foot and leg fit into the boot.\nIn view of the foregoing, the main object of this invention is to provide an articulated pneumomassage sleeve adapted to fit onto and conform to a foot and lower leg of a patient suffering from an excess of body fluid, the sleeve functioning to displace excess fluid and relieve the patient\"\"s condition.\nAmong the significant advantages of an articulated sleeve in accordance with the invention are the following:\nA. The articulated sleeve can be fitted onto the foot and lower leg of patients whose foot and leg dimensions lie within broad ranges. Hence the articulated sleeve functions as a universal pneumomassage appliance.\nB. The articulated sleeve can be quickly installed on the foot and leg of a patient or quickly removed therefrom without any difficulty for the patient such as the difficulty experienced when trying to put on a combination of a cylindrical sleeve and a boot.\nC. When the articulated sleeve is mounted on the foot and leg of a patient, the foot section is then hinged to the leg section, so that the patient is then free to flex his foot in the course of treatment.\nD. The articulated sleeve is relatively inexpensive to manufacture, for incorporated in its structure are standard zippers or other fasteners, and the sleeve is made of bow cost materials.\nMore particularly, an object of this invention is to provide an articulated sleeve having on its rear side a row of normally-fastened transverse articulation joints, articulation of the sleeve being effected by unfastening a selected joint.\nAlso an object of this invention is to provide a pair of pants having an upper section that extends to the waist of the patient and a pair of legs below the upper section, each formed by an articulated sleeve enveloping inflatable cells.\nBriefly stated, the objects are accomplished by a pneumomassage articulated sleeve adapted to fit onto a foot and lower leg of a patent afflicted with lymphoedema or other disorder resulting in, excess body fluid. Enveloped by the sleeve is a series of overlapping inflatable cells which when the sleeve is worn, are sequentially inflated to create massaging forces giving rise to a sequential or peristaltic massage action pumping the excess fluid away from the foot and leg.\nAlong the rear of the sleeve is a row of transverse slots forming articulation joints, each of which is normally fastened. To accommodate the sleeve to the patient to be treated a slot in the row is unfastened to define a foot section conforming to the foot of the patient and a leg section hinged to the foot section conforming to the lower leg of the patient.\nA sleeve in accordance with the invention can be used for other pneumomassage purposes, as for example, for cosmetic massaging."}
{"text": "This specification relates to online content provisioning.\nThe Internet provides access to a wide variety of resources, such as video and/or audio files, and web pages for particular subjects. Access to these resources has provided opportunities for content items to be provided with the resources. There are many different types of content items that can be provided, such as video items, text items, audio items, image items, and combinations thereof. For example, some online content items include a text portion and an image portion. The performance of an online content item is often affected by the form of media used to present the content item.\nSome advertisements only include a textual portion. In such cases the content item could be enhanced by accompanying the textual portion by an image portion to display jointly to users. However, unless a content item provider specifies an image or images to be provided with the textual portion, selecting a suitable image to display is a time-intensive process that involves a manual search of images and review of the images and the textual portion."}
{"text": "Diabetes mellitus is a disorder in which the pancreas cannot create sufficient insulin (Type I or insulin dependent) and/or in which insulin is not effective (Type 2 or non-insulin dependent). In the diabetic state, the victim suffers from high blood sugar, which can cause an array of physiological derangements associated with the deterioration of small blood vessels, for example, kidney failure, skin ulcers, or bleeding into the vitreous of the eye. A hypoglycemic reaction (low blood sugar) can be induced by an inadvertent overdose of insulin, or after a normal dose of insulin or glucose-lowering agent accompanied by extraordinary exercise or insufficient food intake.\nConventionally, a person with diabetes carries a self-monitoring blood glucose (SMBG) monitor, which typically requires uncomfortable finger pricking methods. Due to the lack of comfort and convenience, a person with diabetes normally only measures his or her glucose levels two to four times per day. Unfortunately, such time intervals are so far spread apart that the person with diabetes likely finds out too late of a hyperglycemic or hypoglycemic condition, sometimes incurring dangerous side effects. It is not only unlikely that a person with diabetes will take a timely SMBG value, it is also likely that he or she will not know if his or her blood glucose value is going up (higher) or down (lower) based on conventional method. This inhibits the ability to make educated insulin therapy decisions."}
{"text": "This invention relates to pencil sharpeners, particularly to manually operated pencil sharpeners.\nIn the prior art there have been known manually operated pencil sharpeners which are provided with a receptacle for receiving shaving scraps which result from the use of the sharpener. One such prior art pencil sharpener with a receptacle and dust-proof cover is described in U.S. Pat. No. 2,525,854 to Becker. The Becker pencil sharpener has a relatively short, large diameter cylindrical configuration, and it does not lend itself easily to placement within a shirt or jacket pocket in conjunction with pencils and pens.\nIt is an object of the present invention to provide a manually operated pencil sharpener which can be conveniently carried around in a breast pocket along with pens, pencils and other similarly shaped writing implements.\nIt is a further object of the inventon to provide such a pencil sharpener having interchangeable tools for sharpening different writing implements."}
{"text": "1. Field of the Invention\nThe present invention relates to method for the manufacture or treatment of a material web, in particular of a paper or cardboard web, while using a respective machine, in particular a paper machine. The present invention further relates to an apparatus for measuring the hardness and/or the compactness of at least one clothing, in particular, of at least one felt, fabric, belt and/or the like, of a machine, in particular, of a paper machine, for the manufacture or treatment of a material web, in particular of a paper or cardboard web.\n2. Description of the Related Art\nPrevious technologies for studying machine clothing involved online studies of moisture and/or permeability of machine clothing in order to determine their condition and performance. These studies were unable to provide information regarding the compaction (or density) of the machine clothing in cross direction or in machine direction over time. Apart from this, such measurements were unable to provide any indication of nip profile in situations where the clothing had passed through a pressing operation. A method of this kind is described, for example, in U.S. Pat. No. 5,611,893 and DE 691 29 788 T2."}
{"text": "1. Field of the Invention\nA system and method for storing and sequencing a plurality of luggage items. More specifically a system and method for automatically storing, sequencing, and retrieving luggage items from a storage area.\n2. Description of the Prior Art\nWhen passengers travel via commercial conveyance methods such as on airliners, the conveyance company often allows the passengers to check their luggage items. The conveyance company and their agents, collect, process, transport, and return to the passenger the checked luggage items upon arrival at the final passenger's final destination.\nDuring processing of the luggage items by the conveyance company, it is often necessary to hold luggage items in process for a period of time. An example of this occurs when a passenger checks in very early for a flight. That passenger's luggage items must be held until the airline is ready to load the flight.\nIn the current state of the art, stored luggage is managed a number of ways. For example, the stored luggage items may be manually stored in an area designated for this purpose. Such manual handling methods tend to be labor intensive, prone to worker injury, and prone to manual miss-handling error.\nAccording to another example method, as illustrated in FIG. 1, stored luggage items may be sorted to, and accumulated on, long sections of belt conveyor. Each section of belt conveyor may hold a group of luggage items that have common attributes such as flight number or departure time window.\nCertain issues associated with such conveyor storage systems have been recognized. For example, such systems require substantial floor space and capital investment. Furthermore, the handling of stored bags in large groups obviates the possibility of processing individual bags or smaller groups of bags. Such handling may be necessary when a passenger changes itinerary or when a flight schedule changes unexpectedly. When these events happen, the entire group, or entire belt section, must be re-sorted to separate the desired luggage item(s) from the group. These handling constraints impair the flexibility of these systems.\nIn view of the foregoing, there remains a need for improvements to systems and methods for storing and sequencing luggage items."}
{"text": "Many industrial, commercial and consumer processes involve the injection of a gas into a liquid. Intimate contact of the gas and liquid is generally beneficial to the process, and the degree to which this can be achieved will depend on the system's ability to overcome the tendency of the gas to separate from the liquid and coalesce into large bubbles or pockets. Intimate contact, i.e., an increase in the interfacial area can be achieved in a variety of ways, such as the use of specially designed nozzles for the injection of the gas, the use of propellers or stirrers for mechanized mixing, or the insertion of specially engineered conduit sections which function as static mixers.\nWhile any of these methods will suffice to shear gas into small bubbles, certain effects are preferably avoided. Pressure drops, such as those inherent in nozzles, should be avoided or minimized since they consume energy and impede flow. Clogging or fouling of flow passages should also be avoided, since this can cause a rapid rise in the pressure drop if not stop up the flow completely. Still further, the possibility of corrosion must be considered, particularly when this results in frequent or time-consuming down-time for disassembly and replacement of parts."}
{"text": "1. Field of the Invention\nThe invention relates to a polished semiconductor wafer for the fabrication of electronic components, having a front surface and a back surface and an edge which forms the periphery of the semiconductor wafer and is part of a profiled boundary of the semiconductor wafer. The semiconductor wafer has a polished front surface, into which the components are formed. Strict demands are imposed on the flatness of the front surface, and these demands are extraordinarily high if it is intended to accommodate electronic structures with line widths of 0.1 μm or below (≦0.1 μm technology). To enable the maximum number of circuits of this type to be integrated, it is necessary to ensure the required flatness to as close as possible to the edge of the front surface.\n2. The Prior Art\nMost efforts to increase the flatness of the side faces of the semiconductor wafer in general and the front surface in particular have consistently been concentrated on substeps involved in the production of a semiconductor wafer which influence the flatness. These substeps include in particular steps such as the lapping and/or grinding and polishing of one or both side faces. At least one polishing step, carried out as a single-side or double-side polishing step, is virtually always carried out. As is made clear by European Patent No. EP 1119031 A2, however, it is also possible for substeps such as the etching of the side faces to have an effect on the flatness, in particular on the flatness in the edge region of the side faces. A semiconductor wafer is usually etched prior to a first polishing step, in order to remove damage from the surface left behind by a previous shaping operation, for example by grinding and/or lapping of the semiconductor wafer. The abovementioned patent application states that a raised portion in the edge region of the polished front surface of the semiconductor wafer is likely if the semiconductor wafer is exposed to a flow of a liquid etchant which is guided onto the boundary of the semiconductor wafer during etching. To avoid this effect, a shield can be positioned in front of the boundary of the semiconductor wafer preventing the etchant from being able to strike the boundary of the semiconductor wafer directly. The document does not give any indications with regard to any potential for targeted influencing of the flatness of a semiconductor wafer in its edge region, in particular with a view to making the semiconductor wafer suitable for the ≦0.1 μm technology."}
{"text": "1. Field of the Invention\nThis invention relates to circuitry for use with a guidance system for an automatically guided vehicle (AGV) and more specifically to circuitry for detecting the presence of a valid guidepath signal from various types of sensors, such as photocell arrays, infrared light, sound, and the like which operates over a large dynamic range without the need to make gain or threshold adjustments.\n2. Description of the Prior Art\nAutomatically guided vehicles are generally known in the art. Examples of such vehicles are disclosed in the following U.S. Pat. Nos.: 4,328,545; 4,345,662; 3,379,497; 4,623,032; 4,602,334; 4,627,511; 3,039,554; 3,610,363; 3,933,099; 4,003,445; 4,151,526; 4,602,334 and 4,500,970. Such automatically guided vehicles are used in a wide variety of applications, including the transfer of raw materials and subcomponent parts in manufacturing and assembly facilities, the cleaning of floors in warehouses and parking lots and the delivery of mail in business offices. AGVs are also used in a wide variety of applications in the agricultural industry such as, plowing, harvesting, mowing, and the like.\nIn each application the AGV is guided by on-board sensors which follow a guidepath. Various types of systems are used for automatically guiding the vehicle along a guidepath.\nIn U.S. Pat. No. 4,003,445, the guidepath consists of a fluorescent material, applied to a floor, carpet or the like which emits visible light in a predetermined frequency range. The fluorescent material, however, is normally invisible under ambient lighting conditions. In this type of guidance system, an ultraviolet light located on the vehicle, irradiates the guidepath. This radiation stimulates the fluorescent materials in the guidepath causing it to emit visible radiation which is sensed by sensors onboard the AGV.\nOptical guidance systems for use along a predetermined guidepath are susceptible to spurious operation due to background radiation. Thus, the guidance system must be able to detect a \"no-line\" situation. Also such systems are affected by the non-uniform intensity of the fluorescent guidepath. The solution to this problem has heretofore been attempted. For example, U.S. Pat. No. 4,003,445 discloses an automatic gain control (AGC) circuit for distinguishing between varying intensities of the reflected light from the fluorescent guidepath and background radiation. In this system, AGC feedback circuitry is disclosed for continuously adjusting the gain of the sensor circuit to compensate for the variations in the guidepath intensity. The AGC circuitry is also coupled to other circuitry for detecting a no-line situation. However, such circuitry is relatively complicated and requires a continuous gain adjustment of the sensor circuit to compensate for the varying intensity of the guidepath."}
{"text": "1. Field of the Invention\nThis invention relates to computer-assisted methods for identifying and designing small molecule mimetics of erythropoietin.\n2. Description of Related Art\nErythropoietin (EPO) is the primary regulator of the proliferation and differentiation of immature erythroid cells. EPO is produced in the fetal liver and in the adult kidney in response to hypoxia (low oxygen levels in blood or tissue). It circulates in the blood stream where it targets the EPO receptor (EPOR) on committed progenitor cells in the bone marrow and other hematopoietic tissues. Recombinant human erythropoietin (rHuEPO) is widely used in therapy of patients with anaemia due to chronic renal failure, cancer chemotherapy and AZT treatment.\nThe EPO receptor belongs to the cytokine receptor superfamily which includes receptors for other hematopoietic growth factors such as interleukins (ELs), colony stimulating factors (CSFs) as well as growth hormone prolactin and ciliary neurotrophic factor (CNTF). The structural architecture of this family of receptors consists of three modules: a ligand binding extracellular domain, a short trans membrane region and a large cytoplasmic domain. It has been proposed that the extracellular domain of this superfamily comprises two discrete domains each containing approximately 100 residues that fold into a sandwich consisting of 7 antiparallel .beta.-strands with the topology of an Ig constant domain. Members of the family share two characteristic motifs in their extracellular domain: a pair of conserved disulfide bridges in the N-terminal domain, and a WSXWS box (SEQ ID NO:2) (where X is any amino acid residue) in the C-terminal domain. For most members of this receptor superfamily, oligomerization of one or more polypeptide chains is essential for forming high affinity receptor complexes. A homodimer complex has been demonstrated to be the active form of hGHR and a similar model has been suggested for G-CSF, prolactin and EPO receptors.\nErythropoietin induces dimerization of two EPO receptor molecules, which results in subsequent phosphorylation of the cytoplasmic domains by the association with two tyrosine kinase (JAK2) molecules to initiate a cascade of events that leads to the relevant biological.\nGiven the importance of erythropoietin, it would be very desirable to be able to identify molecules capable of binding the EPO receptor and eliciting the response normally elicited by EPO."}
{"text": "The present invention relates to controlling surface defects, as described hereinafter, in steel strip that has a corrosion-resistant metal coating that is formed on the strip by hot-dip coating the strip in a molten bath of coating metal.\nThe present invention relates particularly but not exclusively to metal coated steel strip that can be cold formed (e.g., by roll forming) into an end-use product, such as roofing products.\nThe present invention relates particularly but not exclusively to metal coated steel strip having an aluminum-zinc-silicon alloy coating that can be cold formed (e.g., by roll forming) into an end-use product, such as roofing products. The applicant is interested particularly in aluminum-zinc-silicon alloy coated steel strip that is sold in Australia under the registered trade mark Zincalume and in other countries under the registered trade mark Galvalume.\nThe present invention also relates particularly but not exclusively to metal coated steel strip having an aluminum-zinc-silicon alloy coating with small spangle size, i.e., a coating with an average spangle size of the order of less than 0.5 mm. Coated steel strip products with larger spangle size do not tend to show the generally small defects because the defects are camouflaged by the appearance of the spangle pattern.\nThe term “aluminum-zinc-silicon alloy” is understood herein to mean alloys comprising the following ranges in weight percent of the elements aluminum, zinc and silicon:\nAluminum: 50-60\nZinc: 37-46\nSilicon: 1.2-2.3\nThe term “aluminum-zinc-silicon” alloy is also understood herein to mean alloys that may or may not contain other elements, such as, by way of example, any one or more of iron, vanadium, chromium, and magnesium.\nIn the conventional hot-dip metal coating method, steel strip generally passes through one or more heat treatment furnaces and thereafter into and through a bath of molten coating metal, such as aluminum-zinc-silicon alloy, held in a coating pot. The furnaces may be arranged so that the strip travels horizontally through the furnaces. The furnaces may also be arranged so that the strip travels vertically through the furnaces and passes around a series of upper and lower guide rollers. The coating metal is usually maintained molten in the coating pot by the use of heating inductors. The strip usually exits the heat treatment furnaces via an outlet end section in the form of an elongated furnace exit chute or snout that dips into the bath. Within the bath, the strip passes around one or more sink rolls and is taken upwardly out of the bath. After leaving the coating bath, the strip passes through a coating thickness control station, such as a gas knife or gas wiping station, at which its coated surfaces are subjected to jets of wiping gas to control the thickness of the coating. The coated strip then passes through a cooling section and is subjected to forced cooling. The cooled strip may thereafter be optionally conditioned by passing the coated strip successively through a skin pass rolling section (also known as a temper rolling section) and a tension leveling section. The conditioned strip is coiled at a coiling station.\nThe present invention is concerned particularly but not exclusively with minimizing the presence of particular surface defects on steel strip that has been hot dip coated with an aluminum-zinc-silicon alloy.\nThe particular surface defects are described by the applicant as “rough coating” and “pinhole-uncoated” defects. Typically, a “rough coating” defect is a region that has a substantial variation in coating over a 1 mm length of strip, with the thickness varying between 10 micron thick and 40 micron thick. Typically, a “pinhole-uncoated” defect is a very small region (<0.5 mm in diameter) that is coated."}
{"text": "A reflective mask substrate to be used in a lithography technology using light in a soft X-ray region or vacuum ultraviolet region having a wavelength of from approximately 0.2 to 100 nm (The substrate, lithography and light are hereinafter referred to as “substrate for EUVL”, “EUVL” and “EUV light”, respectively.) is required to have a high flatness of a surface of the substrate such that a strain is not generated on a transferred image even upon irradiation of EUV light and also to have a small surface roughness so that a high reflectance is obtained, namely it is required to be highly flat and highly smooth. In addition, in order to suppress influences due to the temperature change, it is required that a coefficient of thermal expansion is low. For that reason, it is preferably used as the substrate for EUVL a substrate resulting from a quartz glass substrate containing SiO2 as a major component, in particular a quartz glass with low thermal expansion which has been made to have low thermal expansion by incorporating TiO2 as a dopant, through processing to give high flatness and high smoothness. Though the shape, size, thickness, etc. of the substrate for EUVL are not limited, usually used substrate are plate-shaped bodies having a rectangular or square planar shape, a substrate size of typically 152 mm in square, and an effective area of the substrate, which is to be highly flat and highly smooth, is 142 mm in square.\nIn mechanical polishing, even if it is intended to process the surface of the foregoing substrate so as to give high flatness and high smoothness, there is a problem in flatness that a concave-convex shape which is a periodical series of a convex and a concave is formed within the surface of the substrate. In such a concave-convex shape, the distance from one convex to an adjacent convex is typically from 75 to 150 mm. For that reason, it was difficult only by the mechanical polishing to obtain a flatness of not more than 0.15 μm in terms of PV (peak-to-valley), which is required for a substrate of a reflective mask for EUVL. As a method of obtaining a desired highly flat surface free from such a concave-convex shape, Patent Documents 1, 2 and 3 disclose methods of obtaining a flat surface shape by carrying out processing through optimizing local processing amounts while employing a processing method such as gas cluster ion beam etching, plasma etching or polishing with a magnetic viscous fluid.\nHowever, when dry etching is applied to a glass substrate, defects inherent to the dry etching are caused; and when polishing with a magnetic viscous fluid is applied, defects inherent to the polishing with a magnetic viscous fluid are caused, resulting in a problem that the surface roughness becomes large. Patent Document 3 discloses a method of improving such surface roughness or defects by employing non-contact polishing such as EEM.\nPatent Document 1: JP-A-8-293483\nPatent Document 2: JP-A-2002-316835\nPatent Document 3: JP-A-2004-291209"}
{"text": "Directed damage of the skin is used to stimulate regrowth of collagen and to improve skin appearance. A well known method of directed damage is ablating the epidermis using laser radiation having wavelengths strongly absorbed by water so as to heat the water to above boiling temperature. Typical lasers used for epidermis ablation are CO.sub.2 and Er:YAG lasers. Ablating the epidermis using RF (radiofrequency) current is described in U.S. Pat. No. 6,309,387. This treatment significantly reduces wrinkles and improves the skin appearance. The main disadvantages of skin resurfacing are the long healing period that can be over a month long and the high risk of dischromia. These disadvantages have reduced the popularity of ablative skin resurfacing in recent years.\nNon-ablative skin resurfacing is based on heating of the dermis to a sub-necrotic temperature with simultaneous cooling of the skin surface. U.S. Pat. No. 5,810,801 describes penetrating the dermis with infrared laser radiation with dynamic cooling of the skin surface using a cryogen spray.\nWrinkles are created in skin due to the breakage of collagen fibers and to the penetration of fat into the dermal structure. Thus, destroying adipose cells and structure, can improve the surface structure. However, most wrinkle treatment methods target the collagen and do not have a significant effect on deep wrinkles Radio frequency (RF) energy has been used for the treatment of the epidermal and dermal layers of the skin. For example, U.S. Pat. No. 6,749,626 describes use of RF for collagen formation in dermis. This patent describes a method for collagen scar formation. U.S. Pat. Nos. 6,470,216, 6,438,424, 6,430,446, and 6,461,378 disclose methods and apparatuses for affecting the collagen matrix using RF with special electrode structures together with cooling and smoothing of the skin surface. U.S. Pat. Nos. 6,453,202, 6,405,090, 6,381,497, 6,311,090, 5,871,524, and 6,452,912 describe methods and apparatuses for delivering RF energy to the skin using a membrane structure. U.S. Pat. Nos. 6,453,202 and 6,425,912 describe methods and apparatuses for delivering RF energy and creating a reverse temperature gradient on the skin surface. Although a non-ablative treatment is much safer and does not scar the skin tissue, the results of non-ablative treatments are less satisfactory.\nA method described in U.S. patent application No. 20030216719 attempts to maintain the efficiency of ablative treatment with a shorter healing time and a lower risk of adverse effects. The device described in that patent coagulates discrete regions of the skin where the regions have a diameter of tens of micrometers and the distance between the regions is larger than the regions themselves. This treatment provides skin healing within a few days but the results are very superficial and less spectacular than with CO.sub.2 laser treatment, even after multiple treatments.\nU.S. Pat. No. 6,277,116 describes a method of applying electromagnetic energy to the skin through an array of electrodes and delivery electrolyte using a microporous pad.\nA device for ablation of the skin stratum corneum using RF electrodes is described in U.S. Pat. Nos. 6,711,435, 6,708,060, 6,611,706, and 6,597,946. However, the parameters of this device are optimized for the ablation of the stratum corneum so as to enhance drug penetration into the skin, and not for thermal collagen remodeling."}
{"text": "Chemical vapor deposition methods are employed to form films of material on substrates such as wafers or other surfaces during the manufacture or processing of semiconductors. In chemical vapor deposition, a chemical vapor deposition precursor, also known as a chemical vapor deposition chemical compound, is decomposed thermally, chemically, photochemically or by plasma activation, to form a thin film having a desired composition. For instance, a vapor phase chemical vapor deposition precursor can be contacted with a substrate that is heated to a temperature higher than the decomposition temperature of the precursor, to form a metal or metal oxide film on the substrate. Preferably, chemical vapor deposition precursors are volatile, heat decomposable and capable of producing uniform films under chemical vapor deposition conditions.\nThe semiconductor industry is currently considering the use of thin films of various metals for a variety of applications. Many organometallic complexes have been evaluated as potential precursors for the formation of these thin films. A need exists in the industry for developing new compounds and for exploring their potential as chemical vapor deposition precursors for film depositions.\nFor the chemical vapor deposition of silicon-containing films (e.g., SiO2), compounds such as silane, chlorinated silanes, and alkoxy silanes (e.g., TEOS) are well known. However, as next generation oxide materials with higher dielectric constants, so called ‘high-k’ materials (e.g., HfO2), are integrated, and concurrently new precursors are developed for these materials (e.g., hafnium amides), other silicon precursors will require development for the deposition of ternary systems and beyond (e.g., hafnium silicates).\nFor silicon amide compounds with cyclic amide ligands, an example reported in the literature is tetrakis(pyrrolidinyl)silane (a solid at room temperature, mp=30° C.). Inorg. Nucl. Chem. Letters 1969 5 733 discloses tetrakis(pyrrolidinyl)silane compound and a low yield synthetic method for preparation thereof.\nU.S. Patent Application Publication Nos. U.S. 2002/0187644 A1 and U.S. 2002/0175393 A1 disclose metalloamide precursor compositions having stated utility for forming dielectric thin films such as gate dielectric, high dielectric constant metal oxides, and ferroelectric metal oxides and to a low temperature chemical vapor deposition process for deposition of such dielectric thin films utilizing the compositions.\nIn developing methods for forming thin films by chemical vapor deposition or atomic layer deposition methods, a need continues to exist for precursors that preferably are liquid at room temperature, have adequate vapor pressure, have appropriate thermal stability (i.e., for chemical vapor deposition will decompose on the heated substrate but not during delivery, and for atomic layer deposition will not decompose thermally but will react when exposed to co-reactant), can form uniform films, and will leave behind very little, if any, undesired impurities (e.g., halides, carbon, etc.). Therefore, a need continues to exist for developing new compounds and for exploring their potential as chemical vapor or atomic layer deposition precursors for film depositions. It would therefore be desirable in the art to provide a precursor that possesses some, or preferably all, of the above characteristics."}
{"text": "Integrated lens arrays or aspherical lenses are used in various projectors as important optical components to condense light from a high-brightness light source or to diffuse such light uniformly. Nowadays, such lenses for optical components are required to have high precision. In particular, microlens arrays utilized for image processing are required to have high optical precision to provide high-definition images.\nPrecision press-molding, wherein polishing is not necessary, has been widely adopted as a method for producing such optical lenses. In order to produce high-precision lenses by this method, it is necessary to decrease the thermal deformation of molds and to enhance transferability. Therefore, an important technical problem is to lower the deformation temperature of the glass material.\nLow-melting point borosilicate crown glasses are known as glass materials having low deformation temperatures. However, since the chemical durability of low-melting point borosilicate glasses tends to decrease with lowering deformation temperature, they are not suitable for practical use in lens arrays or aspherical lenses.\nIn contrast, conventional glass materials such as Pyrex (registered trademark) and white plate glasses have good chemical durability. However, such glass materials have high deformation temperatures; therefore, when they are used for press molding, mold thermal deformation tends to occur, which makes it difficult to produce high-precision lenses.\nJapanese Unexamined Patent Publication No. 1993-19379 discloses, as a glass for press molding, an optical glass for precision press molding comprising SiO2, B2O3, and Al2O3 as principal components, and further comprising 3% to 12% by weight of Li2O. Although this glass is described as having a deformation temperature of 560° C. or lower and allowing precision press molding at a temperature of 600° C. or lower, an improvement is desired in terms of press moldability at low temperature. This glass is also unsatisfactory in terms of chemical durability.\nJapanese Unexamined Patent Publication No. 2001-89183 discloses an optical glass for press molding comprising, as an essential component, an alkali metal oxide comprising Li2O, Na2O, and K2O. Although this glass has improved press moldability at low temperature, a further improvement is desired in terms of chemical durability."}
{"text": "Without limiting the scope of the invention, its background is described in connection with the development of genetically modified whole cell cancer vaccines. More specifically, the present invention relates to vaccines capable of augmenting tumor antigen expression, presentation, and processing through expression of the GM-CSF transgene and attenuating secretory immunosuppressive activity of TGF-β via furin bi-functional shRNA transgene induced knockdown.\nThe prevailing hypothesis for immune tolerance to cancer vaccines include the low immunogenicity of the tumor cells, the lack of appropriate presentation by professional antigen presenting cells, immune selection of antigen-loss variants, tumor induced immunosuppression, and tumor induced privileged site. Whole cancer cell vaccines can potentially solicit broad-based, polyvalent immune responses to both defined and undefined tumor antigens, thereby addressing the possibility of tumor resistance through downregulation and/or selection for antigen-loss variants. A method for making a master cell bank of whole cell vaccines for the treatment of cancer can be found in U.S. Pat. No. 7,763,461 issued to Link et al. (2010). According to the '461 patent, tumor cells are engineered to express an α (1,3) galactosyl epitope through ex-vivo gene therapy protocols. The cells are then irradiated or otherwise killed and administered to a patient. The α-galactosyl epitope causes opsonization of the tumor cell enhancing uptake of the opsonized tumor cell by antigen presenting cells which results in enhanced tumor specific antigen presentation. The animal's immune system thus is stimulated to produce tumor specific cytotoxic cells and antibodies which will attack and kill tumor cells present in the animal.\nGranulocyte-macrophage colony-stimulating factor, often abbreviated to GM-CSF, is a protein secreted by macrophages, T cells, mast cells, endothelial cells and fibroblasts. When integrated as a cytokine transgene, GM-CSF enhances presentation of cancer vaccine peptides, tumor cell lysates, or whole tumor cells from either autologous patient tumor cells or established allogeneic tumor cell lines. GM-CSF induces the differentiation of hematopoietic precursors and attracts them to the site of vaccination. GM-CSF also functions as an adjuvant for dendritic cell maturation and activation processes. However, GM-CSF-mediated immunosensitization can be suppressed by different isoforms of transforming growth factor beta (TGF-β) produced and/or secreted by tumor cells. The TGF-β family of multifunctional proteins possesses well known immunosuppressive activities. The three known TGF-β ligands (TGF-β1, -β2, and -β3) are ubiquitous in human cancers. TGF-β overexpression correlates with tumor progression and poor prognosis. Elevated TGF-β levels within the tumor microenvironment are linked to an anergic tumor response. TGF-β directly and indirectly inhibits GM-CSF induced maturation of dendritic cells and their expression of MHC class II and co-stimulatory molecules. This negative impact of TGF-β on GM-CSF-mediated immune activation supports the rationale of depleting TGF-β secretion in GM-CSF-based cancer cell vaccines.\nAll mature isoforms of TGF-β require furin-mediated limited proteolytic cleavage for proper activity. Furin, a calcium-dependent serine endoprotease, is a member of the subtilisin-like proprotein convertase family. Furin is best known for the functional activation of TGF-β with corresponding immunoregulatory ramifications. Apart from the previously described immunosuppressive activities of tumor secreted TGF-β, conditional deletion of endogenously expressed furin in T lymphocytes has been found to allow for normal T-cell development, but impaired function of regulatory and effector T cells, which produced less TGF-β1. Furin expression by T cells appears to be indispensable in maintaining peripheral tolerance, which is due, at least in part, to its non-redundant, essential function in regulating TGF-β1 production.\nHigh levels of furin have been demonstrated in virtually all cancer lines. The inventors and others have found that up to a 10-fold higher level of TGF-β1 may be produced by human colorectal, lung cancer, and melanoma cells, and likely impact the immune tolerance state by a higher magnitude. The presence of furin in tumor cells likely contributes significantly to the maintenance of tumor directed, TGF-β-mediated peripheral immune tolerance. Hence furin knockdown represents a novel and attractive approach for optimizing GM-CSF-mediated immunosensitization and vaccine development. Chen et al. (2004) in U.S. patent application Ser. No. 20040242518 provide methods and compositions for inhibiting influenza infection and/or replication based on the phenomenon of RNAi as well as systems for identifying effective siRNAs and shRNAs for inhibiting influenza virus and systems for studying influenza virus infective mechanisms. The invention also provides methods and compositions for inhibiting infection, pathogenicity and/or replication of other infectious agents, particularly those that infect cells that are directly accessible from outside the body, e.g., skin cells or mucosal cells. In addition, the invention provides compositions comprising an RNAi-inducing entity, e.g., an siRNA, shRNA, or RNAi-inducing vector targeted to an influenza virus transcript and any of a variety of delivery agents. The invention further includes methods of use of the compositions for treatment of influenza.\nInterferon-gamma (γIFN) is a key immunoregulatory cytokine that plays a critical role in the host innate and adaptive immune response and in tumor control. Also known as type II interferon, γIFN is a single-copy gene whose expression is regulated at multiple levels. γIFN coordinates a diverse array of cellular programs through transcriptional regulation of immunologically relevant genes. Initially, it was believed that CD4+ T helper cell type 1 (Th1) lymphocytes, CD8+ cytotoxic lymphocytes, and NK cells exclusively produced γIFN. However, there is now evidence that other cells, such as B cells, NKT cells, and professional antigen-presenting cells (APCs) secrete γIFN. γIFN production by professional APCs [monocyte/macrophage, dendritic cells (DCs)] acting locally may be important in cell self-activation and activation of nearby cells. γIFN secretion by NK cells and possibly professional APCs is likely to be important in early host defense against infection, whereas T lymphocytes become the major source of γIFN in the adaptive immune response. Furthermore, a role for γIFN in preventing development of primary and transplanted tumors has been identified. γIFN production is controlled by cytokines secreted by APCs, most notably interleukins (IL) IL-12 and IL-18. Negative regulators of γIFN production include IL-4, IL-10, glucocorticoids, and TGF-β."}
{"text": "The present invention relates to a method for processing received signals transmitted via a transmission channel and also relates to a corresponding device. In particular, the present invention is directed to a method and corresponding device for channel impulse response improvements in TDMA systems.\nRecently, mobile radio telecommunication systems have widely spread. Such mobile radio telecommunication system operate for example according to a commonly agreed standard, like for example the GSM standard. According to GSM standard, data transmission is performed according to a method of time divisional multiple access (TDMA). The TDMA transmission principle specifies that data are transmitted from a transmitter to a receiver and vice versa only during respectively specified time slots of frames.\nData transmission in such telecommunication systems substantially relies on digital data transmission. However, between a mobile radio transceiver device as a subscriber terminal (hereinafter: mobile station MS) and a stationary radio transceiver device as a serving radio network element (hereinafter: base station BS) data have necessarily to be transmitted as analogue data via the air interface Um.\nConsequently, data transmitted by the mobile station MS are received via a base station antenna means of a reception device of the base station BS as analogue data. In the course of the further processing of the thus received data by the reception device, the analogue data are analog to digital converted, i.e. passed through an A/D converter means. At the next stage of the processing, the obtained digital data are supplied to an equalizer means for being equalized. The thus obtained output data are then supplied to a channel codec means for coding/decoding the data. At the reception side, decoding is performed to separate received speech data from associated signaling data.\nParticular attention in the course of this processing has to be paid to the equalizing of the received data, since the equalizing is required to reconstruct, at the reception side, the transmitted signal from a mixture of received signals.\nFor example, assuming a situation in a radio telecommunication network with a base station BS and only a single mobile station MS present in the radio coverage area of the base station. Then, a signal s transmitted from the mobile station MS may reach the base station BS directly via line of sight at a time s(t). However, the same signal s may be deflected by, e.g., a building, a mountain or the like present in the environment. Thus, the same signal may reach the base station BS at a later point of time s(t+T), and is thus superposed to the signal s(t). Due to the delay T, both received signals are no longer in phase with each other. Moreover, the delayed signal s(t+T) may even be more attenuated than the signal s(t) due to the longer transmission path. Thus, the signal received by the base station BS and originating from the mobile station MS is distorted. Now, assuming that another mobile station MSxe2x80x2 is additionally present, then signals sxe2x80x2(txe2x80x2), sxe2x80x2(txe2x80x2+Txe2x80x2) are additionally received by the base station BS, which may lead to interference between the respective transmitted data symbols (intersymbol interference).\nTherefore, an equalizer means has to reconstruct (detect) the initially transmitted signal s(t) and/or sxe2x80x2(txe2x80x2) from the received mixture of signals s(t), s(t+T), sxe2x80x2(txe2x80x2), sxe2x80x2(txe2x80x2+Txe2x80x2).\nThe thus reconstructed (or detected) signal is required to be as similar to the originally transmitted signal as possible. This reconstruction is therefore a main concern when designing equalizers, e.g. for use in a reception device of a base station BS.\nHitherto, in equalizers of reception devices used in mobile telecommunication systems, the channel impulse response (CIR) is estimated, and the thus estimated channel impulse response is used to detect, i.e. to reconstruct the transmitted data symbols.\nDocument WO 94 28661 representing the prior art described in the preambles of claims 1 and 7 discloses a method for estimating a channel state based on a low complexity model. This method is used by a received in Rayleigh fading environments, and the receiver contains at least one channel stated estimator which uses an auto regressive model for channel variations.\nDocument U.S. Pat. No. 5,329,547 discloses a method and an apparatus for coherent communication. According to this method, reference symbols are inserted into a data stream. The channel response is estimated by utilizing the stream of reference samples. An estimated data symbol is detected from the stream of data samples by utilizing the estimated channel response.\nDocument WO 97 23089 discloses a method and apparatus for channel identification utilizing two Least-Squares (LS) estimators. Each LS estimator is used for calculating a sequence of channel values, and further for determining and estimated channel impulse response, over an entire frequency band thereof.\nThe estimated channel impulse response is usually based on the received samples of the stream of data symbols. Therefore, it is an estimate of the actual (observed) channel impulse response. However, an error in the channel impulse response leads to a degraded performance when detecting/reconstructing transmitted data symbols from received data symbols.\nIn a previous solution, the above described estimated channel impulse response is assumed to be ideal, thereby disregarding errors, and it is used as it is for data symbol reconstruction. However, this solution was unsatisfactory due to the error in the observed channel impulse response being not taken into account.\nAnother known solution for improving the performance of an equalizer is disclosed in document U.S. Pat. No. 5,251,233 by Labedz et. al., assigned to Motorola Inc. The basic idea described therein is to estimate the energy of taps in the impulse response. Those taps having an energy below a predetermined threshold level are zeroed in order to reduce noise in the estimated channel impulse response. This previously proposed method may be applied to complex taps or real taps, i.e. may be performed separately for real and imaginary taps. The method disclosed in the Labedz patent improves (i.e. reduces) the bit error rate BER in the reproduced detected signal (which is an indication for the quality of a receiver) in situations where the channel interference ratio C/I or C/(I+N), respectively, with N being a noise signal superposed to an interfering signal, is low and the channel has a poor quality. The method is also beneficial in channels having a short impulse response. (The term taps refers to coefficients of each respective delay element of the corresponding FIR model for the used circuitry. Thus, based on the tap values, poles and zeroes of the transmission function may be calculated.)\nHowever, under good channel conditions, the method as proposed by Labedz et. al. even degrades the receiver performance which may be attributable to the zeroing of amplitude coefficients of certain taps.\nConsequently, it is an object of the present invention to provide a method for processing received signals transmitted via a transmission channel, and to provide a corresponding device which improve the receiver performance in all channels irrespective of actual channel conditions.\nAccording to the present invention, this object is achieved by a method for processing received signals transmitted via a transmission channel, the method comprising the step of obtaining an observed response function of said transmission channel, based on said received signals; characterized by the steps of deriving an estimated variance of said received signals; and modifying said observed response function by applying said estimated variance to said observed response function, thereby obtaining a modified impulse response; wherein said step of modifying is effected by multiplying said estimated response function with a vector, the components of which being determined on the basis of parameters of the estimated response function and said estimated variance.\nFurthermore, the present invention proposes a device for processing received signals transmitted via a transmission channel, said device comprising obtaining means adapted to obtain an observed response function of said transmission channel, based on said received signals; characterized by deriving means adapted to derive an estimated variance of said received signals; and modifying means adapted to modify said observed response function by applying said estimated variance to said observed response function, thereby obtaining a modified impulse response; said modifying means further comprises a multiplying means adapted to multiply said estimated response function with a vector, the components of which being determined on the basis of parameters of the estimated response function of said estimated variance.\nAdvantageous further developments of the present invention are as set out in the respective dependent claims.\nAccording to the present invention, the proposed method and/or device are very easy to implement. Moreover, the present invention improves the receiver performance in all channel types as specified in the GSM 5.05 recommendation (e.g. channel models xe2x80x9cHT100xe2x80x9d and xe2x80x9cRA250xe2x80x9d). Additionally, the improvement may be achieved not only independent of the channel type as such, but also independent of the specific channel condition, i.e. in situations of different channel interference ratios.\nPreferred embodiments of the present invention are described herein below in detail by way of example with reference to the accompanying drawings."}
{"text": "There are three important subpopulations of T Lymphocytes (T cells): helper T cells, which interact with B cells to amplify production of antibody; effector T cells, which carry out the direct cell-killing function of T cells and make certain lymphokines (non-antibody products) which are responsible for delayed hypersensitivity; and supressor T cells, which participate in the regulation of both antibody-medicated and cell-mediated immunity.\nT cells must be activated before any of these forms of activity are expressed. Usually the activation follows from exposure to antigen, but other less specific factors such a interleukin 2 are also believed to participate in the activation of T cells.\nUsually there is a latent period of around a week to 10 days after first exposure to antigen before the T cells develop initial reactivity. Shortly after that the reactivity subsides. Upon a second contact with antigen, T cells show an accelerated memory response with high activity developing within 2-5 days.\nThe addition of antigen to cultured lymphocytes induces a small proportion of T cells to differentiate into the large rapidly dividing blast cells. T cells can also be transformed by culturing them together with the lymphocytes of individuals of the same species, a so-called mixed lymphocyte culture. Because of extensive polymorphism at HLA loci, the two cell populations are virtually always different antigenically and they stimulate each other to undergo blast transformation. For example, isolated blood lymphocytes from recipient and prospective donor are maintained together for several days in tissue culture. Blast transformation occurs if allogenic cells are present. This process may be referred to as alloactivation. Mixed lymphocyte cultures may be established by treating one set of cells in a manner that prevents blast transformation of that set of cells such as irradiating the cells.\nA T lymphocyte will recognize an antigen only if the antigen is properly presented by a presenting cell which in many cases is a macrophage. The antigen must be presented juxtaposed to a compatible Ia molecule, a surface molecule coded for by one of the class I transplantation or histocompatibility genes. In man, there are at least two distinct families of Ia molecules encoded by mixed lymphocyte culture genes, HLA-DR and MT or DS. These genes control the formation of the specialized complementary Ia structures on the surface of a presenting cell and the T cells that provide for proper presentation of antigens. T cells may interact with B cells, or other T cells, if the cell possesses complementary Ia structures and if it recognizes the same antigenic determinant or a different determinant on the same antigenic molecule.\nDuring the process of activation, T cells develop new surface antigens, so-called T cell activation antigens. Most of these T cell activation antigens, however, are not T cell specific. For example, the transferrin receptor, the insulin receptor and the 4F2 antigen appear on proliferating cells of many types. Only one T cell activation antigen, Tac (Interleukin-2 receptor), is found only on activated T cells.\nThe kinetics of antigen appearance on activated T cells has been studied in order to gain some insight into the function of these molecules. Cotner et al. examined the kinetics of appearance of several T cell antigens and found that each exhibited a characteristic and reproducible time of appearance (9). Based upon this, the antigens could be classified as early, intermediate or late appearing antigens.\nThe relevancy of various cell surface markers to the functional heterogeneity of alloreactive T cell clones has been examined and the rearrangement of the T B-chain gene of the antigen receptor in wild and mutant variants of clones with altered phenotypes and function compared (7). BALB/C mice have been immunized with helper T cell clones to produce two monoclonal antibodies which detect distinct cell surface molecules associated with T cell helper activity and with the production of T replacing factor (8).\nThe 4F2 antigen, Tac, the 49.9 antigen and the transferrin receptor appear within 24 hours of mitogen stimulation, before the onset of DNA synthesis, and are classified as early antigens. Early appearing antigens may be associated with cell growth. The HLA-DR antigen and the 19.2 antigen (Ia antigen) do not appear until about 72 hours after activation and are classified as late antigens. Expression of the OKT 10 antigen by activated T cells is intermediate.\nThe molecule or molecules associated with T cell helper function are unknown. Human helper T cells are defined as lymphocytes which express the T4 surface antigen, a T cell specific glycoprotein molecule expressed by cells which recognize and are restricted by major histocompatibility complex class II antigens. This molecule is also found on killer and suppressor lymphocytes suggesting that the T4 molecule is not involved in helper function. Similarly, the T cell antigen receptor does not seem to be involved in any specific T cell function because the gene which encodes the B-chain of the receptor is rearranged and expressed by helper, suppressor and cytotoxic T cells."}
{"text": "In order to satisfy owners and producers of digital content, such as music, movies, and the like, digital content distribution requires many commercial and technical safeguards. Providers of digital content often must take measures to prevent access to their content by unauthorized computer programs, for fear that these applications will illegally copy, display, or distribute the protected content. Moreover, legitimate content providers often license content from content owners and/or producers (e.g., artists, record labels, movie studios, etc.), and the terms of the licenses typically require that adequate safeguards be used so that the licensed content is protected from pirating. Thus, distributors of such content often use digital rights management (“DRM”) tools to prevent copying and recording of their content. But DRM techniques and the strict terms of the licenses between content providers and content owners/producers often impede legitimate content providers from offering fast, convenient, and ubiquitous access to their content catalogs. The ideas disclosed herein help alleviate these problems, and allow legitimate content providers to expand the ways in which they provide licensed digital content to users and third parties."}
{"text": "For several years, there has been a growing interest, in the field of cosmetics, for depigmenting products: age spots (lentigo senile), chloasma, freckles, pigmentation spots that appear on mature skins after a prolonged exposure to the sun, are badly tolerated by affected people.\nFashion phenomena favouring light complexions in Asia and Africa are also factors in favour of the development of depigmenting products.\nMelanin is a pigment synthesized by melanocytes which are specific cells of the epidermal basal layer.\nHyperpigmentations are often caused by hyperactivity of the melanocytes, eventually associated to an increase of their number.\nWithin melanocytes, oxidation reactions that lead to the formation of melanin are mainly catalyzed by tyrosinase. Resulting pigments are uniformly distributed over the epidermis surface, except when the mechanism is disrupted: melanin is then anarchically accumulated in some areas.\nThe depigmenting properties of some substances have been known for long: as soon as 1936, Oettel noticed a decrease of pigmentation among black cats hair which had been given hydroquinone (Hemsworth B. N. , J.Soc Cosmet. Chem., 1973, 24, 727-733).\nSome substances showed a depigmenting action but cannot be used in cosmetic because of their toxicity (for example, mercury salts), of the cutaneous irritation they induce (mercaptoamines, oxidizing products such as hydrogen peroxide) or because of the excessive slowness of their action (ascorbic acid or ascorbic acid derivatives).\nOther compounds such as phenolic derivatives, corticosteroids, analogs of vitamin A, some vegetal extracts are used as components in many depigmenting preparations now on the market.\nHowever, these compounds show some side-effects: among phenolic derivatives, hydroquinone, which is an inhibitor of tyrosinase, induces quite a few secondary reactions (inflammatory reactions, burns, prickling sensations, irreversible dyschromia in spots); corticosteroids induce such secondary reactions and cannot be used in cosmetic; analogs of vitamin A (beta carotene, canthaxanthine) are able to block the melanocyte activity but the problem with these substances is making them to penetrate into the cells without being metabolized by the cell enzymes. Plant extracts, containing either arbutin (glucose hydroquinone) or flavonoids or also cinnamic derivatives, are also used as depigmenting agents but have a very slow action.\nGenerally speaking, the problem with these depigmenting agents of the prior art is that it will take at least a two months treatment before seeing the first results in vivo."}
{"text": "This invention relates to a structure to surround the fusion plasma of a fusion reactor, or other fusion plasma containing device, said structure commonly called a blanket.\nFusion reactors of all varieties produce energetic neutrons which can advantageously be captured in a blanket region or structure substantially or completely surrounding the reactor core. In a fusion-fission hybrid reactor, the blanket contains a fertile fuel intended to breed fissile fuel and to produce energy by neutron induced fission. In a pure fusion reactor, the blanket contains fertile species which capture neutrons to form valuable isotopes. In both reactor types, the blanket is cooled by a coolant by which means heat is transferred away to cool the reactor, and perhaps generate useful power.\nWhile many nuclear transformations are possibly of interest in a fusion reactor blanket, the generation of tritium from a neutron-induced nuclear reaction of lithium nuclei is especially important since tritium is a fusion reactor fuel.\nThe design of a suitable blanket is considered a significant obstacle to the development of a practical fusion power reactor. Several concepts have been proposed encompassing liquid lithium or solid lithium compounds for tritium breeding, solid fertile fuel for fissile element breeding, and gas, liquid, and even pebble bed coolants. All have technical handicaps. A solid blanket must be replaceable for isotope recovery and for blanket repair which, it develops, is a difficult task. A liquid lithium blanket can be continually or intermittently processed for tritium recovery and can readily be used to recover heat deposited therein by conventional means. However, liquid lithium and other liquid metals in the presence of the strong magnetic fields (in magnetic confinement fusion reactors) experience magnetohydrodynamic forces which limit the serviceability of that type of blanket. Gas cooled blankets have intrinsically lower material densities which reduce the efficiencies of heat transfer. Pebble bed blankets generally require high coolant pumping power.\nConsequently, it is desired to provide a blanket for a fusion reactor which can be refueled on-line, has high density, and is well adapted for power and isotope production."}
{"text": "In recent years, electronic presentations have been made using multiple display devices (e.g., monitors, etc.). In such cases, an operator giving the presentation connects an external display device, such as an external monitor or projector for an electronic device, such as a laptop personal computer (PC) carried by the operator and causes an external output device to display the screen content of the laptop PC by operating the laptop PC in, for example, a mirror display mode.\nIn some conventional electronic devices, such as computer equipment, an operating system for a multi-monitor system can be used to connect two or more display devices to one terminal so that the display devices can be used on the same system in case one display device is separated, which can be intended by a user. For example, Japanese Patent Application Publication No. 2004-62177 describes a multi-monitor display in which a small-screen display device is disposed on a side or upper surface of a large-screen display device so that the display screens of these display devices are driven by the same operating system.\nIn a case in which a multi-monitor setting of the laptop PC is the mirror display mode, for example, the screen content of the laptop PC can be unintentionally displayed on the external display device immediately after the external display device is connected to the laptop PC. Since screen contents of the laptop PC may include content that the operator does not want to be displayed on the external display device, the operator generally wants to quickly switch out of the multi-display mode. This switching operation, however, usually involves pressing a special key or the like, and is not currently known to users. Operation of a special key, in some cases, involves pressing an operation function key “Fn” as a special key and an F-number key (e.g., “F7 ”) at the same time, and this pressing causes the multi-display mode to switch OFF.\nEven when the multi-display mode switches OFF by operating the special key, if the mirror display mode has been selected, the display content of the laptop PC can be displayed on the external display device at least for a moment, and a screen content not intended by the operator is displayed on the external display device."}
{"text": "In recent years, valve opening/closing timing control devices that make it possible to change the opening/closing timing of an intake valve and an exhaust valve in accordance with the operation status of an internal combustion engine (hereinafter referred to as an “engine” as well) have been put to practical use. These valve opening/closing timing control devices have a mechanism that, for example, by changing the relative rotation phase of the driven rotating body with respect to the rotation of the driving rotating body (hereinafter referred to as simply “relative rotation phase”) by means of an engine operation, changes the opening/closing timing of an intake/exhaust valve that is opened and closed accompanying the rotation of the driven rotating body.\nIn general, the optimal opening/closing timing of the intake/exhaust valve differs according to the operation state of the engine, such as the state in which the engine is started, and the state in which the vehicle is traveling. By constraining the relative rotation phase to an intermediate lock phase between the maximum retard phase and the maximum advance phase when starting the engine, the opening/closing timing of the intake/exhaust valve that is optimal for starting the engine is realized, and a case in which a knocking sound is generated due to a partition of a fluid pressure chamber formed by the driving rotating body and the driven rotating body swinging is suppressed. For this reason, it is desired that the relative rotation phase is constrained to the intermediate lock phase before the engine is stopped. The stopping of the engine also includes idling stops in which the engine is stopped for a short amount of time at an intersection or the like so as to suppress the discharge of exhaust gas or the consumption of gasoline.\nPatent Document 1 discloses a valve timing adjustment device that can reliably perform locking when a lock pin is to be locked in an intermediate phase between the maximum advance phase and the maximum retard phase. With this valve timing adjustment device, a control valve is configured to connect an advancing port and a lock port to a main supply port and a discharge opening respectively by moving to a first region, and to connect both the advancing port and the lock port to the main supply port by moving to a second region that is shifted away from the first region in a second direction. The first region is a lock region for locking the rotation phase in a restricted phase using a first main restricting member. Furthermore, the first region has a reduction region in which an advance supply flow amount to be supplied to an advancing chamber by connecting the advancing port in communication with the advancing chamber to the main supply port is reduced to a flow amount that is smaller than the flow amount at the moving end in the first direction. Accordingly, in the constricted region, the speed at which a vane rotor rotates to the advance side is a slower speed that corresponds to the flow amount controlled to be a smaller amount. Furthermore, when the phase of the vane rotor gradually changes to the advance side in this way, the lock port and the discharge opening are connected so that working oil is discharged from the lock chamber. Accordingly, locking of the rotation phase that accompanies the flow of the working oil from the lock chamber can be performed reliably due to the phase of the vane rotor gradually changing to the advance side."}
{"text": "Augmented reality (AR) technology augments an image of a real-world environment (reality) by superimposing supplemental information (such as, pictures, videos, three-dimensional (3D) models, and other sensory enhancements) onto the real-word environment. The AR technology overlays virtual objects onto the image of the real world, enhancing a user's perception of reality and providing a user with an immersive, interactive experience."}
{"text": "For many years, radiation therapy has been an indispensable part of cancer therapy. Especially, X-rays and Gamma rays, consisting of small particles called photons, have been used for attacking the tumor. They remove electrons out of tumor cell atoms and destroy larger molecules within the cells by destroying smaller chemical compounds before. Radiation beams also damage the genetic code, the DNA of the cells and thereby the construction plans for essential proteins. Furthermore, cells cannot replicate any more and consequently, they die. However, the energy is only partly transferred to the tumor. This problem has been partly solved by modern techniques in which photon beams hit the tumor from multiple directions and meet at a defined target where they discharge a maximum of energy. At the same time mobile apertures screen the sensitive healthy tissue from radiation. The so called Intensity Modified Radiation Therapy (IMRT) improves the treatment results of conventional radiotherapy considerably.\nIon radiation does not use photons, but positively charged ions, atomic nuclei which have lost at least one electron from the atomic shell. The particles mainly used are hydrogen atomic nuclei (protons) and carbon atomic nuclei, which are very heavy. This particular type of ions is therefore called heavy ion. Atomic nuclei are accelerated in large devices to about three quarters of the speed of light and shot into the tumor. The depth of penetration can be enhanced by speeding up the ions. Ion beams have always been interesting candidates for radiation therapy, since they have special physical characteristics: When they hit the body they travel very fast through the outer layers and lose hardly any energy before they decelerate in the depth and eventually get stuck and transfer their entire deleterious energy to the surrounding tissue. Scientists call this moment the Bragg peak after its discoverer, the English physicist William Henry Bragg. Therefore, ion beams are well-suited for treating tumors located deeply inside the body. Also, tumors with irregular edges can be scanned accurately to the millimeter with the Intensity-Controlled Raster Scan Method.\nRadiation therapy with protons and carbon ion beams has been shown to be an effective treatment for tumors. In addition, such therapy has been shown to result in less damage to surrounding healthy tissue than conventional gamma radiation therapy.\nRadiation planning for determining the magnitude and position of a radiation dose to be administered is typically based on previous MRI imaging or CT imaging, which may have taken place at a considerable time period before the therapy takes place. In the intervening period, the tissue to be irradiated may have moved, or changed shape. This may result in the irradiation of healthy tissue and/or missing diseased tissue, which may prevent the disease from moving to a remission.\nIn the past several years, techniques have been developed for more precisely targeting incident radiation upon tumors in a human body. Such techniques have been achieved by using advances in X-ray sources and collimation systems. Also, imaging of the body has been improved by optimizing the related targeting algorithms by mapping the body with CT systems and inputting the data into the algorithm(s). Recently, MRI systems and X-ray systems have been presented that combine attributes of both types of imaging.\nUnfortunately, for the treatment of tumors such as those from malignant cancers, X-ray therapy does not possess the precision of dose administration or stopping power that particle-based therapy has. It is difficult or impossible to irradiate targeted tissue and transfer the energy of the radiation to the desired tissue in a precise manner, as X-rays by their nature as highly energetic photons, pass through most soft tissue unimpeded.\nWhile radiations systems and methods do exist to irradiate tumors and other tissue with particles, e.g., protons, carbon nuclei, etc., such charged particles by definition will naturally deflect from a straight line trajectory in the presence of a high-strength magnetic field such as produced by common MRI systems. Such deflection can render ion beam targeting inaccurate, and can potentially cause the irradiation of healthy tissue while at the same time diseased tissue can escape radiation exposure."}
{"text": "This Application claims the benefit of the Japanese Patent No. 10-109785, filed Apr. 20, 1998, the entire contents of which are hereby incorporated herein by reference.\nThe present invention relates to a vector quantization method used to quantize a linear predictive coefficient in speech encoding and, more particular, to a vector quantization method having constraints in quantization vectors.\nA technique for performing linear prediction analysis of speech, decomposing the analysis result into a residual signal and a linear predictive coefficient representing a spectrum envelope, and processing them has been popular. A CELP (Code Exited Linear Prediction) scheme extensively studied recently in the fields of speech encoding is also based on the linear prediction analysis. A linear predictive coefficient and residual signal are quantized by VQ (Vector Quantization). In the CELP scheme, the linear predictive coefficient is often transformed into an LSP (Line Spectrum Pair) parameter, and then quantized. According to one of the reasons for this, stability of a synthesis filter can be easily discriminated.\nIn the CELP scheme, the synthesis filter is arranged on the basis of the LSP parameter on the decoding side. The quantized residual signal is passed through the synthesis filter to generate decoded speech. For this reason, unless the synthesis filter is stable, the decoded speech oscillates to greatly degrade the speech quality.\nAs is well known, if an LSP parameter w={w1, w2, . . . , wp} obtained by linear prediction analysis of degree p satisfies the following condition:\no less than w1 less than w2 less than . . .  less than wp less than xcfx80xe2x80x83xe2x80x83(1)\nthe synthesis filter is stable. The order of magnitudes of components of the LSP parameter represented by condition (1) (to be referred to as an LSP parameter order hereinafter) is checked to allow easily determining the stability of the synthesis filter.\nWhen an interval between the parameter components w1, w2, . . . , wp decreases although the LSP parameter order is not reversed, the synthesis filter abruptly becomes unstable. Care must be taken for quantizing the LSP parameter components having small intervals. Small quantization errors generated by the LSP parameter components having small intervals greatly influence stability of the synthesis filter.\nTo solve the above problem, conventionally, a predetermined value D is defined for an interval between the adjacent LSP parameter components. When the interval between the adjacent components of the quantized LSP parameter is smaller than the predetermined value D, an appropriate countermeasure is made. Japanese Patent No. 2,659,605 (reference 1) describes a method of performing correction processing to increase a small interval between the adjacent components to assure the predetermined value D. This method is simple and easy, but distortion by correction processing is not evaluated in quantization.\nJpn. Pat. Appln. KOKAI Publication No. 6-120841 (reference 2) discloses a technique for solving the above problem. In reference 2, stability check and, as needed, correction processing are performed for each quantized LSP parameter obtained from a codebook. The distance between the corrected quantized LSP parameter and the input LSP parameter is calculated. With this technique, the distortion generated by correction processing can be included in quantization distortion and then evaluated, thereby improving quantization efficiency. According to this method, however, the calculation quantity increases because stability check must be performed for all quantized LSP parameters within the search loop of the codebook.\nAs described above, if a constraint such as the LSP parameter order or the interval between adjacent components is imposed on the quantization vector, it must be checked if the quantization vector satisfies the constraint after quantization or within the search loop of the codebook in order to assure stability of the synthesis filter. This check must desirably be performed in a smaller quantity of calculation. In particular, the calculation quantity increases with an increase in the number of candidates of the codebook in the check within the search loop. A small difference in calculation quantity becomes a large difference as a whole. A reduction in calculation quantity in the check is an important problem. According to the conventional methods, however, it is difficult to check if the quantization vector satisfies the constraint in a small calculation quantity, as described above.\nIn quantizing an LSP parameter upon nonlinear transform such as logarithmic transform, it is difficult to calculate the interval between the adjacent components in the transform range.\nIt is an object of the present invention to provide a vector quantization apparatus capable of efficiently obtaining a quantization vector satisfying a constraint in a smaller calculation quantity and a vector quantization method therefor.\nAccording to the present invention ,there is provided a vector quantization apparatus comprising a codebook which stores a plurality of code vectors; a constraint relieving section which relieves a predetermined constraint imposed on a quantization vector from an input vector input to an input terminal to generate a target vector, a error calculating section which calculates an error between the target vector and a code vector extracted from the codebook, and an error evaluation section for evaluating this error, selecting from the codebook a code vector constituting an unconstrained quantized vector for approximating the target vector, and outputting an index representing the code vector.\nAccording to the invention, there is provided a vector quantization method comprising quantizing an input vector upon appropriately transforming the input vector on an encoding side in advance, and transforming a decoding result in an inverse manner to the transform of the input vector to obtain a quantization vector satisfying a constraint, thereby reducing the calculation quantity as compared to the conventional case.\nMore specifically, on the encoding side, an input vector is transformed with a predetermined transform function to generate a target vector. At least one code vector constituting a first quantization vector that approximates this target vector is selected from a codebook, and an index representing the selected code vector is output. The transform function transforms the first quantization vector in an inverse manner to the transform of the input vector to generate a second quantization vector satisfying a predetermined constraint.\nThe decoding side receives the index representing at least one code vector constituting the first quantization vector that approximates the target vector generated by transforming the input vector with the predetermined transform function. This code vector is extracted from a codebook. The first quantization vector constituting the code vector is transformed in an inverse manner to the transform function to generate the second quantization vector. Note that the transform function is designed to allow the second quantization vector to satisfy the predetermined constraint.\nThe vector quantization method according to the present invention is particularly suitable for LSP parameter vector quantization. In quantizing an LSP parameter vector, an LSP parameter serves as an input vector on the encoding side. A constraint vector representing a predetermined constraint is subtracted from the input vector to generate a target vector. At least one code vector constituting a first quantization vector having a minimum error with respect to the target vector is selected from a codebook. An index representing the selected code vector is output. Note that the constraint vector is designed such that the interval between the adjacent components of the quantized LSP parameter constituting a second quantization vector generated by synthesizing the first quantization vector and the constraint vector has a predetermined value or more.\nThe decoding side receives the index representing at least one code vector constituting the first quantization vector having the minimum error with respect to the target vector obtained upon subtracting the constraint vector representing the predetermined constraint from the input vector made of the LSP parameter. This code vector is extracted from a codebook. The constraint vector is added to the first quantization vector made of this code vector to generate the second quantization vector. Note that the constraint vector is designed such that the interval between the adjacent components of the quantized LSP parameter constituting the second quantization vector is the predetermined value or more.\nThe vector quantization method according to the present invention is also applicable to LPC parameter predictive coding. In this case, on the encoding side, an LSP parameter serves as an input vector. A constraint vector representing a predetermined constraint is subtracted from the input vector to generate a target vector. At the same time, a predictive vector is generated using a vector obtained by subtracting the constraint vector from a previous quantization vector. At least one code vector synthesized with this predictive vector and constituting a first quantization vector having a minimum error with respect to the target vector is selected from a codebook. An index representing this code vector is output. Note that the constraint vector is designed such that the interval between the adjacent components of the quantized LSP parameter constituting the second quantization vector generated by putting together the first quantization vector and the constraint vector is a predetermined value or more in the same manner as described above.\nThe decoding side receives the index representing at least one code vector constituting the first quantization vector having the minimum error with respect to the target vector generated by subtracting the constraint vector representing the predetermined constraint from the input vector made of the LSP parameter. This code vector is extracted from a codebook, and at the same time, the predictive vector is generated using the vector obtained by subtracting the constraint vector from the past quantized vector. The first quantization vector is generated by synthesizing this code vector and the predictive vector. The first quantization vector and the constraint vector are synthesized to generate the second quantization vector. Note that the constraint vector is designed such that the interval between the adjacent components of the quantized LSP parameter constituting the second quantization vector has the predetermined value or more in the same manner as described above.\nAccording to still another vector quantization method of the present invention, an LSP parameter serves as an input vector on the encoding side. A constraint vector representing a predetermined constraint is subtracted from this input vector. The difference is nonlinearly transformed to generate a target vector. The constraint vector is subtracted from a past quantized vector. The difference is nonlinearly transformed to generate a predictive vector. At least one code vector synthesized with this predictive vector and constituting a first quantization vector having a minimum error with respect to the target vector is selected from a codebook. An index representing this code vector is output. Note that the constraint vector is designed such that the interval between the adjacent components of the quantized LSP parameter constituting the second quantization vector generated by putting together the first quantization vector and the constraint vector is a predetermined value or more in the same manner as described above.\nThe decoding side receives the index representing at least one code vector constituting the first quantization vector having the minimum error with respect to the target vector generated by subtracting the constraint vector representing the predetermined constraint from the input vector made of the LSP parameter and nonlinearly transforming the difference. This code vector is extracted from a codebook, and at the same time, the predictive vector is generated using the vector obtained by subtracting the constraint vector from the past quantized vector and nonlinearly transforming the difference. The first quantization vector is generated by putting together this code vector and the predictive vector and then subjected to a nonlinear inverse transformation, and then the first quantization vector and the constraint vector are synthesized to generate the second quantization vector. Note that the constraint vector is designed such that the interval between the adjacent components of the quantized LSP parameter constituting the second quantization vector has the predetermined value or more in the same manner as described above.\nAs described above, according to the present invention, the input vector is transformed to relieve the predetermined constraint from the quantization vector in encoding. In decoding the quantization vector, the quantization vector is decoded and then transformed in an inverse manner to the transform in encoding the input vector to generate the quantization vector satisfying the constraint. Whether the quantization vector satisfies the constraint can be checked in a small calculation quantity.\nMore specifically, when LSP parameter vector quantization is taken into consideration, the constraint for keeping the interval between the adjacent components of the LSP parameter to the predetermined value must be imposed on the quantization LSP vector in addition to the condition of the LSP parameter order in order to assure stability of the synthesis filter. The calculation quantity necessary for checking if the latter condition is satisfied becomes an important problem. According to the present invention, this check can be performed by comparing the magnitudes of the adjacent components, thereby greatly reducing the calculation quantity.\nIn quantizing the LSP parameter upon its nonlinear transform such as logarithmic transform, the constraint is subtracted before the transform, thereby eliminating the calculation for the interval of the transformed LSP parameter in quantization.\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out hereinafter."}
{"text": "The demand for systems that use a variety of devices at a location to monitor a variety of conditions, such as monitoring homes and businesses for alarm conditions, allowing users to centrally control various devices (such as thermostats, switches, cameras, appliances, etc.), monitor medical conditions, and the like has continued to grow as more home and business owners seek better control over their premises and to protect it from various hazards and threats. Such hazards and threats include intrusion, fire, carbon monoxide and flooding, among others dangers that may be monitored and reported to a monitoring station.\nConventional systems typically employ a control panel and/or gateway that receive “event” (such as triggering alarms) and other information from various sensors and devices, and is used to operate them. This may be done locally by the user, or remotely via a monitoring center. In the case of alarm events, the monitoring center can also take appropriate action, such as notifying emergency responders. Installation and servicing complexity associated with these systems tends to be high, as an installer has to physically position, mount, and configure the control panel and all of the various sensors, while taking into account a variety of performance characteristics and requirements for each device to ensure proper operation of the system. These systems also typically incorporate a manufacturer's specific technology designed for the manufacturer's security application, and only certain devices may only appropriately interoperate with other devices in certain ways. This is true as well for more recent all-in-one (AIO) security systems, in which the control panel and a user interface (such as a keypad) are combined in a single unit, even portable AIO systems where the control panel may be relocated around the premises and not permanently installed. For example, such units may sit on top of a table or on the floor, but nevertheless communicates with life safety sensors in a similar manner as a wall-mounted security panel.\nHowever, these existing security systems suffer from the same problem, namely, a lack of health monitoring of the security system. In particular, these existing security systems only generate a system alert when the health levels of sensors or components fall below predefined minimum thresholds. For example, a battery alert for a sensor will only trigger an alert when the battery is falls below an operational threshold, but this operational threshold typically corresponds to a battery level of where the sensor or component is forced to turn off or will turn off shortly. In other words, health level thresholds in existing systems are only triggered when the levels are so bad such that at least one sensor or component of the system stops functioning properly.\nFurther, customers have no way of knowing when the heath levels of their security systems are going to degrade to the point that functionality, i.e., monitoring capabilities, of the security system are affected. Such a failure may come at an inopportune time such as when the customer is on vacation or otherwise away from the premise. Therefore, the user is not able to fix the problem, e.g., change batteries, or call for a service technician because the customer may not even know a problem exist until the customer returns to the premise. Such a sudden failure in monitoring capabilities of these existing security systems may inadvertently provide thieves the window of opportunity they have been waiting for."}
{"text": "Field\nThe present disclosure relates to an extension device for a bone anchor, in particular for use in minimally invasive surgery (MIS). The present disclosure also relates to a system including such an extension device and a bone anchor, where the bone anchor includes an anchoring section and a receiving part for receiving a rod to couple the rod to the anchoring section. The extension device includes a first sleeve and a second sleeve positioned within the first sleeve. The extension device also includes a coupling structure that allows the extension device to couple to the receiving part of the bone anchor and that inhibits translational and rotational movement of the extension device relative to the receiving part.\nDescription of Related Art\nExtension devices, also called head extenders, for pedicle screws for use in minimally invasive surgery are known in the art. For example, U.S. Pat. No. 7,563,264 B2 describes a spinal stabilization system for a minimally invasive procedure wherein detachable sleeves may be coupled to a collar of a bone anchor to allow for formation of the spinal stabilization system through a small skin incision. The detachable sleeve members may allow for alignment of the collars to facilitate insertion of an elongated member in the collars. A coupling system is provided between the sleeve and the collar that inhibits translational movement of the sleeve relative to the collar. In one embodiment, the sleeve may be coupled to a collar of a bone fastener assembly with movable members that may be threaded into threaded openings in the collar.\nWO 2013/112689 A2 describes a minimally invasive tower access device comprising an elongated outer sleeve that slidably receives an elongated inner sleeve. A lock nut is used to secure the inner sleeve and outer sleeve in a locked mode."}
{"text": "A considerable need exists for a dosage form useful for treating cardiovascular diseases. The dosage form should be therapeutically indicated for treating cardiovascular diseases including angina pectoris, hypertension, congestive heart failure, and possess vasodilation properties for decreasing systemic vascular resistance.\nThe beneficial drug isradipine is therapeutically indicated for treating cardiovascular diseases. The cardiovascular diseases include angina pectoris, hypertension and congestive heart failure as disclosed in a patient study reported in The American Journal of Cardiology, Vol. 59, pp 70B-74B, (1987). The drug was administered in the study in a bulk, nonrate uncontrolled dose that was subjected to the changing adverse environment of the gastrointestinal tract.\nIn light of the above presentation it will be appreciated by those versed in the dispensing art to which this invention pertains that a pressing need exists for a rate controlled dosage form that can deliver the valuable drug isradipine to a patient in critical need of cardiovascular therapy. The pressing need exists also for an oral dosage form that can deliver isradipine at a controlled rate in a constant dose per unit time over a prolonged period of time for its beneficial hemodynamic effects substantially independent of the variable environment of the gastrointestinal tract. It will be appreciated further by those versed in the dispensing art that such a novel and unique dosage form that can administer isradipine in a rate controlled dose over time, and simultaneously provide cardiovascular therapy, would represent an advancement and valuable contribution to the art."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for detecting combustion pressure in an internal combustion engine. More particularly, the present invention concerns a combustion pressure detecting method and apparatus for an engine whereby errors in detection may be rectified.\n2. Description of the Prior Art\nThere are various known pressure sensors for detecting combustion pressure in an engine such as those of piezo-electric type, semiconductor type, capacitive type and others. However, every one of these known sensors encounters difficulty in detecting the pressure of combustion with a reasonably high accuracy over an extended period (use life) under adverse environmental conditions such as high voltage, high temperature and high pressure. Thus none of these sensors is particularly suited for installation to motor vehicles. In addition to the aforementioned sensors, there may be mentioned a pressure sensor in which optical fibers are used. By way of example, Japanese Unexamined Patent Publication No. 59-60334 discloses a typical one of the hitherto known pressure sensor in which optical fibers are used.\nThe known pressure sensor employing optical fiber is generally intended for use in process controls. When this sensor is employed as a combustion pressure detecting device in an engine, there arises a problem that temperature of the pressure detecting device itself varies over a very wide range under the influence of combustion temperature in the engine so the device has the disadvantage that it is impossible to accurately measure the combustion pressure."}
{"text": "Along with social progress and technology development, more and more people use access terminals to wirelessly connect to the Internet or enterprise internal networks to acquire information. In order to improve browsing speed and experience of mobile device browsers, pre-reading function has already been achieved by the mobile device browsers, i.e. webpage pre-reading can be carried out by reading a keyword in advance, and after a pre-reading page is completely loaded, a webpage can be switched to enter the pre-reading page, and the contents of the pre-reading page can be displayed.\nHowever, in implementing a webpage pre-reading process on existing mobile device browsers, due to the fact that nearly all the page contents (for example, pictures of products and the like) are loaded in a layzyload mode, JS (JavaScript) cannot load or display all of the contents (for example, part of the product pictures and the like) in the pre-reading webpage. Thus, the user experience is influenced."}
{"text": "High quality printing paper must have a number of physical properties. Two of the most important are a flat and smooth surface to facilitate printing in a press and gloss to produce a more attractive surface, particularly after printing. These properties can be obtained by a variety of techniques, such as coating the paper with pigments and binder and finishing it in one or more pressing operations.\nOne of the most common finishing operations employed in the manufacture of printing paper is supercalendering, in which paper is passed through a series of nips formed by steel rolls pressed against cotton filled rolls at very high pressures, typically at nip loads between 175 KN/M and 437.5 KN/M (1000 and 2500 pounds per lineal inch). This typically results in nip pressures of 13,780 KN/M.sup.2 to 27,560 KN/M.sup.2 (2000 to 4000 p.s.i.).\nTraditional supercalender stacks are not externally heated, but heat is generated when the cotton filled rolls subjected to the extremely high pressures in the nip flex intermittently with each revolution. The nip temperatures in such supercalenders typically reach levels of about 71.degree. C. (160.degree. F.). Another important element in producing good results is having a high moisture content in the paper as it passes through the supercalender. Typically, the moisture content will be 7% to 9%, or higher, of the bone dry fiber weight. Flatness, smoothness and high gloss are obtained in supercalenders because of extreme compression and densification of the sheet. The densification undesirably results in reduced opacity and a blackening effect in overly moist portions.\nSupercalenders commonly consist of a large number of rolls (9 to 14), alternating steel and resilient, in order to obtain the desired smoothness and gloss. In order to obtain smoothness on both sides it is necessary to run an even number of rolls and with two resilient rolls (so called \"cushion rolls\") running together midway in the stack to perform the necessary reversing of the side toward the steel rolls. This action is only partly successful at providing two smooth sides since the first side finished towards the steel is later deformed by the exposure to the resilient rolls.\nBecause of this shortcoming and the inherent mechanical problem associated with the \"cushion roll\" nip, many supercalenders operate today with an uneven number of nips and no \"cushion roll\" nip, which results in only one side being finished against the steel, and while gloss values may be manipulated to be close on the two sides, inevitably one side is noticeably rougher than the other.\nAnother form of finishing is machine calendering wherein the paper web is passed between two normally unheated steel rolls pressed together at high pressures. This process produces smoothness, but little gloss because of the absence of shear in the nip.\nAnother common finishing operation is gloss calendering, which uses heated finishing rolls to produce high gloss finishes on coated paper or board without the high pressure of supercalendering. The nip pressures for commercial machines are typically between about 87.5 to 175 KN/M (500 to 1000 pounds per lineal inch) of nip loading. This typically results in nip pressures of 6,890 KN/M.sup.2 to 13,780 KN/M.sup.2 (1000 to 2000 p.s.i.). The lower pressure causes less densification of the paper, and therefore, better opacity, while the higher temperature softens the coating and permits better gloss enchancement. However, the finishing effect is limited to the coating and the uppermost surface of the web. Thus, the surface of the sheet is not as smooth and flat as that produced in supercalendering and has generally been applied to coated board rather than high quality papers. As a result, gloss calendered sheets do not print as satisfactorily in a printing press as do supercalendered sheets.\nIn recent years, many modifications have been made to gloss calendering, machine calendering and supercalendering operations. Some supercalenders have been heated, primarily to improve the uniformity and control of the temperature. Typically, heated supercalenders reach nip temperatures of about 82.degree. C. (180.degree. F.). Temperatures of some machine calenders or supercalenders have been further increased in an attempt to allow a decrease in pressure to produce the same results. In spite of this modification of supercalendering in the direction of gloss calendering, the fundamental effects of the two processes have remained distinct. Supercalendering uniformly compacts the entire sheet to a high degree, thus flattening the surface fibers and all others, as well as producing gloss on the surface. In contrast, gloss calendering molds, flattens, and glosses the surface of the coating and, in the case of uncoated paper the top surface of the fibrous substrate, but compacts the remainder of the sheet much less than supercalendering.\nExamples of gloss calendering are disclosed in U.S. Pat. Nos. 3,124,504; 3,124,480; 3,124,481; 3,190,212; and 3,254,593. These patents collectively describe apparatus capable of nip temperatures from below the boiling point of water to as high as 232.degree. C. (450.degree. F.) and nip pressures from 1,722 to 17,220 KN/M.sup.2 (250 to 2500 p.s.i.). U.S. Pat. No. 3,124,480 describes finishing steps designed to heat the coating on paper to a temperature which temporarily platicizes at least the surface of the coating in contact with the hot drum. A form of supercalendering in which the rolls are heated to relatively high temperatures is disclosed in U.S. Pat. Nos. 3,442,685 and 3,451,331. These patents disclose a method and apparatus capable of producing high gloss on coated paper by heating at least one roll of a supercalender stack to a temperature between 82.degree. C. and 163.degree. C. (180.degree. F. and 325.degree. F.) to plasticize the coating.\nThe one parameter which has been found to be the most critical in gloss calendering and supercalendering has been the moisture content of the paper. High moisture improves the smoothing and glossing effects of both the coating and the paper substrate. Many developments in supercalendering and gloss calendering involve techniques for increasing the moisture in the web or at least in some portions of it before finishing.\nUnfortunately, moisture is an undesirable control parameter. Small variations in moisture cause large variations in the finished properties of the paper. Also, it is undesirable to have more than about 3.5% to about 4.5% moisture in the finished sheet to avoid uneven reel building and sheet curl from later drying. This amount of moisture is a stable amount, and the sheet will not dry significantly below this level under ambient conditions. To have a finished product with the desired low moisture content and still have the desired high moisture content (e.g. 7% to 9%) to facilitate calendering, many heated calendering operations have increased the drum temperature to dry the moister webs.\nNonuniformity of moisture in the sheet can be even a bigger problem than too much moisture. By nonuniformity, it is meant that the moisture content at one place on the sheet is higher or lower than at other locations across the width of the sheet. The nonuniformity can also exist in the machine direction and the thickness of the sheet. Nonuniformity is most severe when calendering takes place immediately after coating, which is to say when the calender is in line with the coater. If coating is done in a separate operation from calendering, the moisture content of the coated paper has time to equalize thoughout the web before calendering.\nThe above cited U.S. Pat. No. 3,124,504 is primarily concerned with very moist webs (up to 35% or 50% moisture) and includes the concept of drying the web while finishing it. Very high temperatures are employed for drying, but temperatures above the boiling point of water are said to be needed only if the web is wetter than 5% to 8% of the bone dry weight. The web moisture content is also noted as being an important element in the process disclosed in above cited U.S. Pat. Nos. 3,442,685 and 3,451,331. The patents teach that it is best for the paper to have about 7% moisture content, and moisture can be added before the supercalender to improve the finishing effects. The addition of moisture before finishing is also described in above cited U.S. Pat. No. 3,124,481 to manufacture glazed uncoated paper. U.S. Pat. No. 2,214,641 also moistens the surface of the web before finishing. In U.S. Pat. No. 4,012,543, gloss calendering is undertaken immediately after coating before too much of the moisture is lost from the coating. In this disclosure, finishing is carried out at a web moisture content of 9% to 10% of the bone dry weight. In contrast, U.S. Pat. No. 3,268,354, takes special steps to dry the surface of the coating, but to maintain a wet interface between the coating and the fibrous web before gloss calendering. The web in this disclosure has a moisture content of at least 15% at the interface."}
{"text": "FIG. 12 shows a perspective view schematically illustrating an embodiment of a card-type device (refer to Patent Document 1, for example). A card-type device 40 has a wireless communication function, and includes a circuit board 41, a flat housing 42 which accommodates the circuit board 41, and an antenna structure body 43 which is disposed outside of the housing 42 and which performs wireless communication utilizing radio waves. The antenna structure body 43 includes an antenna turning shaft 44 at a base end thereof, and the antenna turning shaft 44 projects from the antenna structure body 43. The housing 42 has a through-hole on one side which allows the antenna turning shaft 44 to be inserted into the housing 42. The housing 42 also has inside thereof a bearing (not shown) which receives the antenna turning shaft 44. With this configuration of the card-type device 40, the antenna structure body 43 turns on the antenna turning shaft 44, and the shaft 44 acts as the axis of turn. For example, when the card-type device 40 is used by being inserted into a card slot of an information terminal apparatus such as a personal computer, the antenna structure body 43 is turned so as to stand on an upper surface of the housing 42 as shown by a dashed line in FIG. 12. Thereby, the sensitivity to radio waves for transmission/reception is enhanced. On the other hand, when the card-type device 40 is detached from the card slot of the information terminal apparatus and is carried about, the antenna structure body 43 is turned into an accommodated flat position to be laid flat as shown by a solid line in FIG. 12. Thus, the card-type device 40 is readily carried.    Patent Document 1: Japanese Unexamined Patent Application Publication No. 2001-69029    Patent Document 2: Japanese Unexamined Patent Application Publication No. 2001-243435    Patent Document 3: Japanese Unexamined Patent Application Publication No. 6-21846\nIn the configuration of the card-type device 40 shown in FIG. 12, in order to turn the antenna structure body 43 from the flat position to the upright position, the antenna structure body 43 is turned upward manually, for example, drawn up using a nail. The manual operation of turning the antenna structure body 43 upright is troublesome, and therefore, the device 40 is not user-friendly."}
{"text": "1. Field\nA method for function monitoring in medical accelerator systems is provided.\n2. Related Art\nAccelerator systems that generate high-energy electromagnetic radiation or particle radiation are used in medicine for radiation therapy. Radiation therapy exposes regions of a patient's body to the radiation. Energy is transmitted to these exposed regions in order to bring about changes in cell components. Radiation therapy is used particularly in treating cancer by killing tumor cells or reducing the capability of the cells to divide. Both linear and cyclic accelerators can be used in accelerator systems.\nConventionally, it has been necessary to manually connect an oscilloscope to a box in the control room of the accelerator system. This connection enables a user to monitor the signals from the accelerator system relevant for the operation. A proper connection requires time-consuming tuning of the oscilloscope. In addition, the signal information obtained by the oscilloscope must be separately recorded, for example, on paper.\nSince the information obtained must be recorded separately, only a small portion of the signal data, for example the amplitude and width, can be picked up or acquired. Generally, this data acquisition is done only at relatively long maintenance intervals, for example, every three months.\nEven with arrangements for transmitting the signals received to a servicing device by modem, intervention by the user is still required. The user must still connect the oscilloscope and forward the data. Data acquisition is limited to the brief times in which an oscilloscope is connected. Because of the difficulty of connecting the oscilloscope and tuning the signal, data acquisition generally only takes place when the accelerator system malfunctions or during routine maintenance at certain relatively long maintenance intervals. Thus, monitoring the accelerator system during the operation of the system is hardly possible because the presence of a system technician in the operation room could bother the medical personnel who are operating the system.\nThus, a method that permits improved function monitoring in medical accelerator systems is desired."}
{"text": "Conventional plug and play technology can be a contradiction in terms. In more detail, plug and play technology can enable creation of interfaces and establishment of communications between multiple devices. Ideally, the creation of the interfaces and subsequent communications occur with a minimal amount of user activity. However, for conventional plug and play technology, this may not always be the case.\nFor example, devices can include or be associated with different components, operating systems, protocols, etc. Thus, a user attempting to interface a personal organizer with a personal computer, for example, may not be able to initiate user interface creation and enable communications between the devices by physically connecting the devices. When interfacing the personal organizer with the personal computer for the first time, the communications protocols and associated software considerations must first be analyzed and configured on both machines before any useful communication between the devices can occur.\nOne drawback of conventional plug and play technology is that software is required for both devices. The software usually accompanies a device in packaging in which the device was purchased (e.g., encapsulated within a compact disk). Inclusion of the software together with the product is associated with various inefficiencies. For instance, it can be costly to create software and package it together with a device. For example, in order to place software together with a device within device packaging, costs associated with making, testing, developing, fixing bugs, customizing, branding, burning, etc. are incurred. Additionally, presenting the software on a storage media is associated with other problems. In order to establish compatibility between two devices, a compact disk or other media that includes the software must be present. This can be inconvenient for a traveler who is attempting to interface two devices for a first time, but did not remember to bring the compact disk that retains the requisite software. Inclusion of software is also inefficient because of increased effort required by the user.\nTo complicate matters, compatibility issues may exist between a device that is desirably interfaced and the compatibility software. For example, the compatibility software can be stored on a CD or DVD drive, whereas a personal computer to which a personal digital assistant is desirably interfaced may not be associated with either drive. Additionally, the personal computer upon which the traveler wishes to install compatibility software may have administrative passwords or some other security functionality that prevents the traveler from using the device. To further complicate matters, it is possible that the traveler lacks the technical knowledge to overcome problems associated with installing compatibility software on a personal computer.\nThe Internet is also a popular medium from which to obtain software that facilitates plug-and-play functionality between devices. The Internet solves some problems associated with conventional plug and play technology, such as not forcing an individual to retain physical possession of the storage media that houses the device compatibility software. Using the Internet, however, as a storage media for compatibility software is not without its faults. One such fault is that at least one of the desirably synchronized devices (e.g., a personal computer) requires an Internet connection in order to obtain the compatibility software. Furthermore, the user may be required to find the appropriate web site in order to locate the compatibility software. Again, the technically deficient user may encounter difficulties in finding the required software in the vastness of the Internet"}
{"text": "1. Field of the Invention\nThe present invention relates to isocyanurate-modified polyisocyanate compositions, to a process for the preparation thereof, and to microcellular foams prepared from said polyisocyanate compositions. More particularly, the invention relates to polyisocyanate compositions containing isocyanurate linkages prepared by reacting said polyisocyanate in the presence of an alkali metal organo-hydroxamate or the metal salt-free acid complex, hereinafter referred to as the \"half-salt.\"\n2. Description of the Prior Art\nThe prior art generally teaches the preparation of isocyanurate-modified isocyanate products, employing a variety of catalysts. U.S. Pat. No. 3,996,223 teaches a process for the preparation of polyurethane quasi-prepolymers containing isocyanurate structures by polymerization in the presence of catalytic quantities of a mixture of Mannich bases and carbamic acid esters from isocyanate and alcohols containing secondary alcohol groups. U.S. Pat. No. 3,645,979 teaches the polymerization of an aromatic polyisocyanate and an aliphatic polyisocyanate in the presence of an organic phosphine catalyst. U.S. Pat. No. 3,652,424 teaches a process for the preparation of polyisocyanate compositions containing 5 to 19 percent by weight of toluene diisocyanate trimer by employing as trimerization catalysts substituted guanidines and isobiguanidines."}
{"text": "With increasing development of science and technology, the processing speed of a computing chip for an electronic device is gradually increased. Consequently, a great amount of waste heat is generated while the computing chip is operated. It is a critical issue to dissipate the heat. For example, a computer host comprises a main board and a display card. The display card is mounted on the main board. Generally, the chip of the display card generates more heat than the main board. Consequently, the heat generated by the chip of the display card has to be dissipated in high priority.\nHowever, for most electronic devices, plural function circuit boards are mounted on the main board in parallel and arranged side by side. Since the airflow for convection is usually obstructed by the adjacent function circuit boards, the heat generated by the chip of the display card cannot be efficiently dissipated away by the existing heat-dissipating mechanism. Therefore, the heat-dissipating mechanism of the electronic device needs to be further improved."}
{"text": "1. Technical Field\nThe present invention relates to a liquid ejection head, such as an ink-jet type recording head, a liquid ejection apparatus, and a manufacturing method of a liquid ejection head, and more particularly to a liquid ejection head which includes a channel unit provided with a series of liquid channels extending from a common liquid chamber down to nozzle orifices through pressure chambers and which can discharge liquid in the form of liquid droplets from the nozzle orifices thereof, a liquid ejection apparatus, and a manufacturing method of a liquid ejection head.\n2. Related Art\nAs for liquid ejection heads structured to discharge out liquid droplets from nozzle orifices thereof by causing pressure fluctuation to liquid accommodated in pressure chambers, there are known, for example, ink-jet type recording heads used for image recording apparatuses, such as printers, color material ejection heads used for manufacture of color filters of liquid crystal displays, electrode material ejection heads used for electrode formation of organic electroluminescence (EL) displays and field luminescence displays (FED), and living body organic matter ejection heads used for manufacture of biochips (biotip), etc.\nThere are various forms of such a liquid ejection head. One exemplary form of the ink-jet type recording head (hereinafter, referred to as a recording head) in an ink-jet type recording apparatus (hereinafter, referred to as a printer) includes a channel unit fixed to a head case, the channel unit having a structure in which a nozzle substrate provided with a plurality of nozzle orifices, a channel substrate provided with a channel section such as pressure chamber cavities or grooves which define a series of ink channels from a common ink chamber down to the nozzle orifices through pressure chambers, and a resilient plate (also called a sealing plate which seals an opening of the channel substrate) which causes elastic deformation to diaphragm sections corresponding to the pressure chambers in response to the operation of a pressure generating means (for example, piezoelectric vibrator) are laminated. Of composition components of the channel unit, the channel substrate is required to be machined at a high density with high precision so as to respond to high density of a recorded image and high speed of recording operation. Accordingly, the channel substrate is suitably made of a crystalline base material such as silicon single crystal base material (silicon wafer) by which a fine shape can be formed by, for example, an anisotropic etching process with sufficient accuracy of dimension.\nAccordingly, it is important to assemble the above-described composition components with high positioning accuracy in order to control the discharge operation of the recording head with high precision. Thus, JP-A-2001-30490 discloses a recording head whose composition components are laminated on and combined with each other in the state in which relative positions of all the composition components are specified in a manner such that composition components of a channel unit and a head case are provided with two through-holes which are positioning references and positioning pins are inserted in the two through-holes.\nBy the way, when a channel substrate is made of a crystalline base material as described above, the through-holes of this channel substrate are formed by etching like pressure chamber cavities. In this case in which there is a difference between the interval of the through-holes of the channel substrate and the interval of the through-holes of the head case, mechanical stress is applied to the channel substrate via the positioning pins due to the difference of the intervals, and thus there was the probability that the channel substrate cracks or breaks."}
{"text": "1. Field of the Invention\nThe present invention relates, in general, to an electric generator using a permanent magnet in a shape of a tapered cone, especially, one which induces a useful equilibrium of magnetic force and, more particularly, to a magnetically balanced electric generator, in which a power-generating core unit and a power-generating magnet are spaced apart at constant intervals using a magnetically repulsive force, which is caused by a permanent magnet shaped like a tapered cone, and a magnetically attractive force, which is caused between the power-generating core unit and a power-generating magnet, and the interval is controlled depending on the intensity of wind, thereby improving power-generation efficiency, the lifetime of parts, and economical efficiency of products.\n2. Description of the Related Art\nGenerally, an electric generator means a device that includes a rotor which rotates under external force, a fixing structure supporting the rotor, and an electricity-generating unit which operates upon rotation of the rotor. However, bearings of such electric generators have a problem in that a certain amount of oil should be regularly supplied to the bearings and the contamination of foreign substances causes damage to the bearings.\nFurther, in case of wind turbines using wind power as an external force, an electrical-connection controller should be installed in order to generate electricity as expressed by the formula E=Cp×(½)×ρ×A×v3, which however is difficult work, and thus problems such as the complexity of parts, burning-out of electric parts or the like occur. Wind turbines are the hardest kinds of electric generators to design because they have to deal with energy proportional to cubic wind velocity.\nMeanwhile, the applicant developed a new construction electric generator in consideration of the fact that existing vertical type wind generators do not sufficiently utilize wind power, which generator was disclosed in Korean Patent Registration No. 0743475.\nFIG. 1 shows the structure of the electric generator in a cross-sectional view.\nThe disclosed electric generator is called a variable electric generator for a wind turbine, which includes a group 11 of magnets which is installed on a lower end of a rotor 10 rotating with the wind, a plurality of core units 21 which is installed in a fixing structure in such a manner as to be arranged concentrically with the magnets of the group and which upon rotation of the rotor 10 interacts with the magnets to thereby generate electricity, an anemometer 30 detecting the magnitude of wind velocity, an electric accumulator 40 storing electricity generated by the core unit 21, and a terminal 50 connecting the respective core units 21 and the electric accumulator together such that the core units are selectively connected with the electric accumulator according to the wind velocity detected by the anemometer 30.\nThe variable electric generator of this type changes a power-generating condition depending on the intensity of wind, having such advantages as more efficiently utilizing wind power.\nIn the meantime, in order to generate electricity using interaction between the power-generating magnet and the core unit, in which a coil is wound around a silicon steel plate, an interval between the power-generating magnet and the core unit should be maintained to approximately 1 mm. Then, since the interval is maintained by a spherical steel bearing 61, which is mounted between the rotor and the fixing structure, the bearing should be periodically lubricated so that the interval can be maintained accurately and efficiently, which is troublesome work, otherwise the bearing may burn out during use, problematically requiring replacement with new one.\nFurthermore, problems such as contact points burning out during manipulation of the power-generating coil, which may result in a reduction in lifetime, may occur."}
{"text": "This invention relates generally to a multi-stage barrel type centrifugal pump and more particularly to a preassembled unit or cartridge forming an inner assembly for a multi-stage barrel type centrifugal pump having all working pump parts including bearings, seals, and thrust balancing device, as well as the impeller and diffuser adaptable for insertion in and removal from the fixed outer casing for the pump.\nIt is well known that in the replacement of the working parts of prior art multi-stage barrel type pumps that it is first necessary to dismantle and remove the bearing seals, balancing device and inner elements, then replace the required working parts, and then reassemble all the elements and associated parts for the pump.\nThe replacement and reassembly process requires adjustments in axial and radial rotor position to obtain satisfactory running clearances between the moving and stationary parts of the pump, and the overall process of element replacement requires many hours of labor in the field where labor rates tend to be high and during which time the pump is not available for service. One effort to meet and overcome this problem is shown in U.S. Pat. No. 3,927,763.\nThe present invention overcomes this problem of barrel type centrifugal pumps by providing a preassembled, adjusted and aligned unit or cartridge having all the working parts of the pump; i.e., bearings, seals, and balancing device as well as the impellers and diffusers; which can be installed or removed easily from the barrel or outer casing for the pump, and renewal of the pump merely requires the removal and installation of a preassembled replacement unit or cartridge.\nThe unit or cartridge concept is made possible by means of a special rig or frame which duplicates the pump casing in supporting the unit or cartridge. This permits all adjustments to be made and checking of the rotor for free rotation to be done before installation in the pump.\nThe fixture can also be moved to the assembly point near the pump and in this context acts as a carrier for transporting the cartridge.\nThus, precise assembly and adjustment can be made at the machine or service shop where more highly trained personnel are located and where labor charges are less costly."}
{"text": "Solid golf balls are typically made with a solid core encased by a cover, both of which can have multiple layers, such as a dual core having a solid center (or inner core) and an outer core layer, or a multi-layer cover having inner and outer cover layers. Generally, golf ball cores and/or centers are constructed with a thermoset rubber, such as a polybutadiene-based composition.\nThermoset polymers, once formed, cannot be reprocessed because the molecular chains are covalently bonded to one another to form a three-dimensional (non-linear) crosslinked network. The physical properties of the uncrosslinked polymer (pre-cure) are dramatically different than the physical properties of the crosslinked polymer (post-cure). For the polymer chains to move, covalent bonds would need to be broken—this is only achieved via degradation of the polymer resulting in dramatic loss of physical properties.\nThermoset rubbers are heated and crosslinked in a variety of processing steps to create a golf ball core having certain desirable characteristics, such as higher or lower compression or hardness, that can impact the spin rate of the ball and/or provide better “feel.” These and other characteristics can be tailored to the needs of golfers of different abilities. Due to the nature of thermoset materials and the heating/curing cycles used to form them into cores, manufacturers can achieve varying properties across the core (i.e., from the core surface to the center of the core). For example, most conventional single core golf ball cores have a ‘hard-to-soft’ hardness gradient from the surface of the core towards the center of the core.\nIn a conventional, polybutadiene-based core, the physical properties of the molded core are highly dependent on the curing cycle (i.e., the time and temperature that the core is subjected to during molding). This time/temperature history, in turn, is inherently variable throughout the core, with the center of the core being exposed to a different time/temperature (i.e., shorter time at a different temperature) than the surface (because of the time it takes to get heat to the center of the core) allowing a property gradient to exist at points between the center and core surface. This physical property gradient is readily measured as a hardness gradient, with a typical range of 5 to 40 Shore C, and more commonly 10 to 30 Shore C, being present in virtually all golf ball cores made from about the year 1970 on.\nThe patent literature contains a number of references that discuss ‘hard-to-soft’ hardness gradients across a thermoset golf ball core. Additionally, a number of patents disclose multilayer thermoset golf ball cores, where each core layer has a different hardness in an attempt to artificially create a hardness ‘gradient’ between core layer and core layer. Because of the melt properties of thermoplastic materials, however, the ability to achieve varied properties across a golf ball core has not been possible.\nUnlike thermoset materials, thermoplastic polymers can be heated and re-formed, repeatedly, with little or no change in physical properties. For example, when at least the crystalline portion of a high molecular weight polymer is softened and/or melted (allowing for flow and formability), then cooled, the initial (pre-melting) and final (post-melting) molecular weights are essentially the same. The structure of thermoplastic polymers are generally linear, or slightly branched, and there is no intermolecular crosslinking or covalent bonding, thereby lending these polymers their thermolabile characteristics. Therefore, with a thermoplastic core, the physical properties pre-molding are effectively the same as the physical properties post-molding. Time/temperature variations have essentially no effect on the physical properties of a thermoplastic polymer.\nAs such, there is a need to achieve a single layer core that has a gradient from the surface to the center, and to achieve a method of producing such a core that is inexpensive and efficient. The gradient may be either soft-to-hard (a “negative” gradient) or hard-to-soft (a “positive” gradient). A core exhibiting such characteristics would allow the golf ball designer to create a thermoplastic core golf ball with unique gradient properties allowing for differences in ball characteristics such as compression, “feel,” and spin."}
{"text": "The present invention relates generally to electronic circuits. More particularly, the present invention provides an adjustable bandwidth phase-locked loop.\nA phase-locked loop (PLL) 10 is an electronic circuit having a phase detector 12, an amplifier 14 and a voltage-controlled oscillator (VCO) 16 generally connected as shown in FIG. 1. The PLL can be used in a wide variety of applications, including: frequency demodulation (both AM and FM signals); frequency synthesis; frequency multiplication; tone decoding and frequency-shift keyed (FSK) decoders. The PLL is a useful building block available from several manufacturers as a single integrated circuit.\nPhase detector 12 provides a comparison between the frequency of input signal 18 and VCO output signal 20, and generates an phase error signal 22 that is a measure of the phase difference between the two input frequencies. Phase error signal 22 may be filtered by a low pass filter 13 then amplified by amplifier 14 resulting in output signal 24. Output signal 24 is Provided as an input to VCO 16 and causes the frequency of VCO output signal 20 to deviate in the direction of the frequency of input signal 18. Filtering of phase error signal 22 smooths swings on the VCO output signal 20. VCO 16 locks on the input frequency maintaining a fixed phase relationship with input signal 18. VCO output signal 18 may be of a different type, e.g., square wave, sine wave or other waveform, than input signal 18. The frequency of VCO output signal 20 is a clean replica of the frequency of input signal 18, which itself may be noisy.\nA PLL circuit may be characterized by two interrelated measurements of the PLL circuit performance: bandwidth and acquisition time. The speed at which the VCO output will lock to the input signal is referred to as the acquisition or settling time for the PLL. A PLL may be configured to operate with a fast settling time by adjusting filter 13. Filter 13 typically includes a series resistor and capacitor coupled between the output of phase detector 12 and ground. Many other filter configurations are used as known in the art. The filter is often referred to as a loop filter. The selection of values for the filter loop capacitors and resistors determines the closed loop bandwidth, and thus the settling time for the PLL.\nBandwidth refers to the range over which the VCO output signal may vary with respect to the input signal from initialization until a lock has occurred. The bandwidth of the PLL circuit has a direct effect on whether lock can be achieved on a given reference signal and how long the acquisition time will be. Bandwidth is directly proportional to the filtering of the phase error signal. The more filtering that is provided, in the form of large RC (or LC) time constants in filter 13, the slower the response of the VCO output signal to variations in the input reference signal. Lock is achieved very slowly or may never be achieved if too much filtering is provided. Conversely, settling times may be improved if less filtering is provided but the output of the VCO may vary and be susceptible to noise in the input signal.\nA wide bandwidth PLL circuit typically provides faster acquisition time while narrow bandwidth systems are slower to achieve lock. Thus, conventional PLL circuits are not well suited for applications requiring both fast settling times as well as narrow bandwidth performance.\nIn general, in one aspect, the invention provides an adjustable bandwidth phase-locked loop including a phase-locked loop having a first input node receiving an input signal having a first frequency, a second input node receiving a feedback signal and an output node which has a signal indicative of an error signal characterizing a frequency error between the input signal and the feedback signal. The adjustable bandwidth phase-locked loop includes a voltage controlled oscillator, coupled to the second input node, receiving the error signal and generating the feedback signal where the feedback signal has a frequency which tracks the first frequency. The adjustable bandwidth phaselocked loop includes a variable loop filter, coupled between the phase-locked loop and the voltage controlled oscillator, filtering the error signal. The variable loop filter is configurable to allow for the tracking of the input signal over both of a broad bandwidth and a narrow bandwidth.\nAspects of the invention include one or more of the following features. The phase-locked loop may include a phase detector which detects differences between frequencies of the input signal and the feedback signal and generates the error signal indicative thereof. The voltage controlled oscillator may include an output signal that tracks the frequency of the input signal. The output signal may be a different form than the input signal.\nThe variable loop filter may include a narrowband loop filter, a wideband loop filter and a switch which switches the narrowband loop filter between the phase-locked loop and the voltage controlled oscillator in accordance with receipt of a bandwidth adjustment signal. The bandwidth adjustment signal may be generated by the phase-locked loop after acquisition of the input signal. A charging circuit may be included to pre-charge the narrowband loop filter to a potential equal to a potential developed across the wideband loop filter so that at a time for switching in the narrowband loop filter the adjustable bandwidth phase-locked loop maintains lock on the input signal. The charging circuit may include an operational amplifier configured as a follower.\nThe narrowband loop filter may include a first capacitor and the wideband loop filter may include a second capacitor. The charging circuit may pre-charge the first capacitor to a voltage level approximately equal to the voltage developed across the second capacitor, where the voltage developed across the second capacitor is in response to the error signal.\nIn another aspect, the invention provides an adjustable bandwidth phase-locked loop including a phase-locked loop having a first input receiving an input signal having a first frequency, a second input for receiving a feedback signal and an output for generating an error signal characterizing a frequency error between the input signal and the feedback signal. The phase-locked loop includes a voltage controlled oscillator receiving the error signal and generating the feedback signal where the feedback signal has a frequency which tracks the first frequency. A wideband loop filter is coupled between the phase-locked loop and the voltage controlled oscillator for filtering the error signal to allow for the tracking of the input signal over a broad bandwidth. A narrowband loop filter is provided for filtering the error signal to allow for the tracking of the input signal over a narrow bandwidth. A controllable switch element is provided, operable to change a mode of the narrowband loop filter between the phase-locked loop and the voltage controlled oscillator wherein the narrowband loop filter effectively minimizes the effect of the wideband loop filter resulting in the tracking of the input signal over a narrow bandwidth.\nIn another aspect, the invention provides a method of adjusting bandwidth of a phase-locked loop circuit where the phase-locked loop circuit includes phase-locked loop for receiving an input signal having a first frequency, a voltage controlled oscillator for generating a locally developed signal having a frequency which tracks the first frequency and a wideband loop filter coupled between the phase-locked loop and the voltage controlled oscillator for dampening the response of the voltage controlled oscillator to allow for the tracking of the input signal over a broad bandwidth. The method includes providing a narrowband loop filter to allow for the tracking of the input signal over a narrow bandwidth and switching in the narrowband loop filter between the phase-locked loop and the voltage controlled oscillator wherein the narrowband loop filter effectively minimizes the effect of the wideband loop filter resulting in the tracking of the input signal over a narrow bandwidth.\nIn another aspect the invention provides a method for quickly locking a narrow bandwidth phase-locked loop circuit where the narrow bandwidth phase-locked loop circuit includes a phase-locked loop for receiving an input signal having a first frequency, a voltage controlled oscillator for generating a locally developed signal having a frequency which tracks the first frequency and a narrowband loop filter coupled between the phase-locked loop and the voltage controlled oscillator for dampening the response of the voltage controlled oscillator to allow for the tracking of the input signal over a narrow bandwidth. The method includes providing a wideband loop filter to allow for the tracking of the input signal over a wide bandwidth and switching in the wideband loop filter between the phase-locked loop and the voltage controlled oscillator while switching out the narrowband loop filter resulting in the tracking of the input signal over a wide bandwidth. Once lock onto the input signal is acquired, the narrowband loop filter is switched back in between the phase-locked loop and the voltage controlled oscillator resulting in the tracking of the input signal over a narrow bandwidth.\nAspects of the invention provide numerous advantages. Applications for the use of PLL devices may be expanded by providing an adjustable bandwidth PLL device. The invention provides for a fast PLL phase or frequency acquisition while maintaining a reduced PLL bandwidth in narrow bandwidth applications. Acquisition time for a narrow bandwidth PLL circuit may be improved by providing an adjustable PLL circuit that switches from wide band to a narrow band performance after acquisition. Other features will be apparent from the following description and claims."}
{"text": "In modern telecommunications, control information and data are typically carried in separate regions of the subframe or packet. For instance, in third generation partnership project (3GPP) long term evolution (LTE) networks, control information is carried in a physical downlink control channel (PDCCH) and/or enhanced PDDCH (ePDCCH) of the subframe, while data/traffic is carried in a physical shared control channel (PSCCH) of the subframe. Additionally, pilot signals are oftentimes communicated contemporaneously with the data transmission in order to facilitate demodulation. More specifically, a pilot signal consists of a series or collection of known reference symbols (i.e., a priori information) which the receiver evaluates upon reception in order to estimate parameters (e.g., fading, scattering) of the air channel. In 3GPP LTE networks, pilot signaling is achieved through the inclusion of cell-specific reference signals (CRS) and/or demodulation reference signals (DMRS) within the subframe.\nData recovery techniques/processes allow receivers to obtain data carried by an earlier data transmission with which the receiver was unable to successfully demodulate/decode. Conventional data recovery techniques rely on either re-transmitting the entire original data transmission or otherwise communicating additional forward error correction (FEC) bits related to the entire original data transmission, e.g., via hybrid automatic repeat request (HARQ) signaling. For instance, one conventional data recovery technique may re-transmit substantially all of the data/traffic carried in the PDSCH region of an LTE subframe. Another conventional data recovery technique may re-transmit FEC bits pertaining to the entire PDSCH region of the LTE subframe. These conventional techniques may consume significant bandwidth, particularly when the original data payload was large or when poor channel conditions necessitate a relatively low coding rate. As such, more efficient mechanisms for recovering data are desired."}
{"text": "The present invention pertains to fiber optical current and voltage sensors and, more particularly, to applications involving filters for use in such sensors.\nOne important aspect of electrical power systems deployed by the power generation industry is the ability to measure power carried over high voltage transmission lines. Power measurement has typically been performed on the high voltage side of the systems before the voltage is stepped down for distribution. However, the need is growing for more frequent and more accurate power measurements (e.g., voltage and current measurements) throughout the power distribution systems. Some recent innovations associated with making such power measurements involve the use of fiber optic sensors as described in more detail in, for example, the above-incorporated by reference U.S. Patent Publication 2011/0204875.\nThere are a number of different electrical power systems which require such monitoring. For example, in nuclear power applications, standby power transformers are employed to run critical equipment in the case of a failure. Normally, these standby power transformers are unloaded, but the high side is connected. By keeping their high sides connected, standby power transformers can be quickly switched on when needed. However, there have been several incidents where these generators were not fully connected at the high side (i.e., they had an open-phase condition), and it was not known until the emergency power was needed. In response to this problem, the Nuclear Regulatory Commission in the United States is implementing a requirement that all such standby generators have monitoring equipment installed to detect the open phase condition.\nIf unloaded, the current flowing into the high side of such standby power transformers is substantially the current attributable to the magnetization current of the power transformer's core, e.g., typically on the order of less than one ampere such as 50 to 800 mA. Alternatively, if the standby power transformer is slightly loaded, e.g., due to the presence of a monitoring load, then the current flowing into the high side of the standby power transformer might be a bit higher, e.g., on the order of an ampere. Regardless, when in the undesirable open-phase condition, the current flowing into the high side of the standby power transformer will drop to approximately zero (there may still be some very small capacitive current).\nThus, in order to tell if a standby power transformer is properly connected or not, a monitoring system can be used to determine when the high side current flow drops below some predetermined threshold, e.g., 25 to 500 mA. In order to reliably measure such small currents, the current sensor must be able to resolve currents several times lower than the low-level threshold, or typically about 10 mA. This value may vary according to the specific transformer under consideration. One way to design an optical current sensor to be able to detect current levels as low as 10 mA is by using a combination of a large number of fiber sensing turns in combination with a number of wire wraps around the sensing fiber loops. For example, Alstom has fielded a product which employs 60 wraps of sensing fiber together with 10 wraps of the primary conductor, yielding a multiplication factor of the current of 600.\nHowever, in some applications it is not practical to employ multiple wraps of the primary conductor around the fiber optic current sensor. Also, it is quite expensive to install a large number of sensing fiber turns (e.g., 600 fiber turns) in a single optical current sensor, which would be needed absent the multiplicative effect provided by the multiple wraps of the primary conductor around the fiber optic current sensor.\nAnother approach to providing an optical current sensor which can detect such low levels of current involves passing the output of the optical current sensor through a comb filter as, for example, described in U.S. patent application Ser. No. 12/598,651, entitled “Adaptive Filters for Fiber Optic Sensors, also referred to herein as U.S. Patent Publication 2011/0204875, the disclosure of which is incorporated here by reference. In U.S. Patent Publication No. 2011/0204875, the comb filter is triggered off the current waveform that itself is being filtered. The raw current passes through a narrow band pass filter around the largest frequency component to be measured (normally the fundamental frequency). The zero crossings of the waveform at the output of this filter are used to track the frequency of the signal, and thus to provide the reset function required in the comb signal processing.\nThe system described in U.S. Patent Publication No. 2011/0204875 is therefore limited in its ability to accurately track the frequency (and thus properly time the resets of the comb filter) when the signal-to-noise ratio becomes too low. At some low level of primary signal, the noise coming through the band pass filter substantially interferes with the zero crossings of the resultant signal, distorting their location, or even adding extraneous zero crossings. Thus the system described in U.S. Patent Publication No. 2011/0204875 suffers a limitation in its ability to recover a small signal from a large noise.\nAccordingly, it would be desirable to provide optical current sensors, transformers, systems and sensing methods which address the foregoing issues."}
{"text": "The human heart is the muscle that is responsible for pumping blood throughout the vascular network. Veins are vessels that carry blood toward the heart while arteries carry blood away from the heart. The human heart consists of two atrial chambers and two ventricular chambers. Atrial chambers receive blood from the body and the ventricular chambers, which include larger muscular walls, pump blood from the heart. A septum separates the left and the right sides of the heart. Movement of the blood is as follows: blood enters the right atrium from either the superior or inferior vena cava and moves into the right ventricle. From the right ventricle, blood is pumped to the lungs via pulmonary arteries to become oxygenated. Once the blood has been oxygenated, the blood returns to the heart by entering the left atrium, via the pulmonary veins, and into the left ventricle. Finally, the blood is pumped from the left ventricle into the aorta and the vascular network.\nFor the vast majority of the population, the events associated with the movement of blood happen without circumstance. However, for many people the heart fails to provide adequate pumping capabilities. These heart failures may include congestive heart failure (commonly referred to as heart disease), which is a condition that results in any structural or functional cardiac disorder that impairs the ability of the heart to fill with or pump blood throughout the body. Presently, there is no known cure for heart disease and long-term treatment is limited to a heart transplant. With only a little over 2,000 patients receiving a heart transplant each year, and over 16,600 more on the waiting list for a heart, there is a persisting need for a cure or at the minimum a means of improving the quality of life of those patients on the waiting list.\nOne such means of bridging the time gap while awaiting a transplant is a circulatory assist system. Circulatory assist systems may also be utilized as a destination therapy for individuals not eligible for a heart transplant. These systems, originally envisioned over thirty years ago, provide assistance to the heart by way of a mechanical pump. In this way, blood is circulated throughout the vascular network despite the diseased heart tissue. Traditionally, these circulatory assist systems include an implantable or extracorporeal pump, a controller (internal or external), and inflow and outflow tubes connecting the pump to the heart and the vascular network. Food and Drug Administration (FDA) approved circulatory assist systems can partially relieve symptoms of breathlessness and fatigue associated with severe heart failure and drastically improve quality of life.\nImplantable pumps may reside in a “pump pocket,” which is a subcutaneous or submuscular space on the chest of a patient, near the clavicle or shoulder. Once the inflow tube is connected to the heart, it normally resides within the thoracic cavity until connected with the pump. In order to access the chest cavity and reach the tube, practitioners must tunnel through the intercostal space adjacent the pump pocket. Once in the thoracic cavity, the practitioner is able to grasp the inflow cannula and direct it to the pump pocket for connection to the pump. When tunneling through the intercostal space, surgeons may resort to general surgical tools at their disposal, such as an anastomotic clamp or forceps. However, such tools may tear or otherwise injure the intercostal tissue when being forced into the intercostal space. Moreover, due to the high amount of force required to tunnel through the intercostal muscle, there is a risk of harming the patient, such as rupturing a vessel. Alternatively, surgeons may use a tunneling device to tunnel through the intercostal space. However, such tools must be able to securely grasp the inflow cannula and direct it to the pump pocket, while still being able to easily disconnect from the inflow cannula so that it may be connected to the pump. Therefore, there is a need in the art for a system and method which allows for a more controlled procedure of connecting a blood pump with the necessary tubing."}
{"text": "A tensioner lever is generally mounted, by means of a mounting bolt or pin, on an engine block or other frame, in the vicinity of a tensioner which cooperates with the lever. The tensioner lever, in cooperation with the tensioner, maintains appropriate tension in the transmission medium to prevent transmission failure due to excess tension and excess loosening of the transmission medium.\nIn one well-known example of a conventional tensioner lever, described in Japanese Laid-open Patent Publication No. 2001-323976 (pages 1 to 4, FIG. 2), a sliding contact portion extending along the traveling plane of the chain, and a reinforcing body, which reinforces the sliding contact portion, are fused together. In another example of a conventional tensioner lever, described in Japanese Utility Model Registration Publication No. 2540896 (pages 1 to 3, FIG. 1), the lever includes a resin shoe for sliding contact with the transmission chain, and an aluminum arm for supporting the resin shoe.\nIn the first of the above two conventional chain levers, the reinforcing body is strengthened by glass fibers. As the reinforcing body continues to pivot about a pivoting shaft on an internal combustion engine, the internal surface of its mounting hole is subject to wear, and the glass fibers are exposed and crushed. The crushed glass fibers act as an abrasive, causing accelerated wear and damage to the mounting hole of the reinforcing body and the pivoting shaft. In the case of a lever incorporating an aluminum arm, the continued pivoting of the lever about a pivoting shaft on an engine causes the arm to be burned and thereby also subjected to wear and damage. The use of a bushing or the like, fitted to the pivoting hole, has been considered to avoid the problem of wear. However, this measure increases the number of parts, the difficulty of assembly, and the production cost. Moreover, in the case of an aluminum arm, recycling of a spent lever is troublesome because the resin shoe, the aluminum arm, and a resin pad, must be separated before disposal.\nAccordingly, the objects of the invention are to solve the above-described problems, and to provide a tensioner lever having excellent mechanical strength and wear resistance, reduced production cost, reduced weight, and ease of recycling. A particular object of the invention is to reduce wear on the inner circumferential surface and the boss portion of the pivoting hole by which the lever is mounted on the engine block."}
{"text": "As optical radar apparatuses capable of setting a detection area in a widthwise direction of a road surface, there are apparatuses such as those disclosed in Japanese Utility Model Laid-Open Nos. 59-117980 and 59-117981.\nIn the apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117980, shown in FIGS. 7A and 7B, light from a light source is condensed to some extent by a lens to obtain a first detection range .theta.T1, and the light source is suitably moved to obtain a second detection range .theta.T2 wider than the first detection range .theta.T1 in a widthwise direction of a road surface by changing the degree of condensing of the lens.\nIn the apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117981, shown in FIG. 8, light from a light source traveling through a lens is diffused by a prism to increase an expansion angle in a road surface direction, thereby setting a detection range wide in a widthwise direction of a road surface.\nIn the apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117980 among the above-described conventional apparatuses, shown in FIGS. 7A and 7B, however, light cannot be emitted simultaneously for the first detection range .theta.T1 and the second detection range .theta.T2, and the detection range is set one-sidedly, since the detection range is changed by moving the light source.\nThe apparatus disclosed in Japanese Utility Model Laid-Open No. 59-117981, shown in FIG. 8, entails the problem of a reduction in the maximum detection distance because light is uniformly diffused by the prism. Because of this problem, it is not possible to meet a demand heretofore made for widening the detection range in a widthwise direction of a road surface in a short-distance area without reducing the maximum detection distance when an optical radar apparatus is used for an inter-vehicle control or an obstacle detecting apparatus.\nThe present invention has been achieved in consideration of the above-described problems, and an object of the present invention is to provide an optical radar apparatus capable of increasing the expansion angle in a road surface direction in a short-distance range without setting the detection range one-sidedly and without sacrificing the maximum detection distance."}
{"text": "Traditionally, computer software is manufactured and distributed through a system where a developer creates software (by designing a solution to a problem and coding that solution), manufactures the software in mass at a production facility, packages the software and then distributes it to the consumer. The software, in the form of computer code, is stored on a secondary storage device, such as a disk. Software products are generally manufactured by copying the computer code to the required number of storage devices. In a normal distribution chain, the consumer, who may be a purchaser or licensee, generally acquires the software (i.e., a storage device containing a copy of the computer code) from wholesale or retail dealers who in turn have acquired the software from the manufacturer.\nA problem with the traditional system for the manufacture and distribution of computer software is that each software developer has its own substantial production, packaging and transportation costs, plus fixed overhead costs for facilities and personnel. When added to the initial investment required to develop software, such costs often make the launching of a new product prohibitive for the developer, especially where the developer is an individual or small corporation. Production and distribution is costly, and these costs are passed on to the consumer by way of higher product prices or fees. Additionally, the developer, where demand for a product is uncertain, may not know how many copies of a piece of software no manufacture and distribute. Production of too many or too few copies results in lost profits.\nOne approach taken by a segment of the software industry in order to reduce the inherent costs associated with the traditional manufacture and distribution of software is to invite consumer participation through the concept of \"shareware\" software. Shareware is software available for copying by consumers free of charge or at nominal costs from electronic bulletin boards accessed using modems or through mail order catalogue services. If, after having had an opportunity to use the software for a short period of time, the user decides to keep the software, the user is then under a good faith obligation to register the software and return a fee to the developer or distributor. In practice, the consumer participation in copying and distributing the available software is good, but the honorarium is rarely paid. As a direct consequence of the poor returns to shareware developers in general, the quality of products and level of support is low.\nAnother and very significant problem facing the manufacturer of computer software is the unauthorized reproduction of software products. Computer software is expensive to design and code, but the duplication of an individual piece of software can be done swiftly and, relative to the initial production costs, at minimal cost. For example, by using operating system commands, a user can simply copy a program from a storage device such as a disk using one of the user's disk drives to another storage device. This illegal copying of software duplicates, on a small scale, a developer's manufacturing process. However, the unauthorized copying of computer software cannot be controlled by the manufacturer and the distribution of authorized copies cannot be traced.\nThe software manufacturing industry is placed in an unfortunate position. The purchase price of software is relatively high to reflect the cost of design and coding. The consumer is often averse to paying a high price for software, especially for software that cannot be evaluated before purchase (because the manufacturer will not allow access to the software before the purchase to prevent pre-purchase copying), and when the consumer is faced with a low risk of being detected and prosecuted for copying software, there is a strong incentive to make an unauthorized copy of the software. In effect, unlicensed, pirated software competes directly for the same market as legally manufactured software. The lost revenue to the manufacturer is passed on to the consumer via higher purchase prices, further increasing the incentive for piracy. With the current trend in the industry to develop software which has an intuitive, graphical interface, not requiring manuals to operate it, this problem can only increase as users do not require access to difficult to copy instruction manuals and tutorials. Further, pirated software can be quickly and widely distributed over telephone lines using modems.\nOne approach the industry has taken to improve the return on investment in the manufacture and distribution of software is to attempt to decrease the amount of software piracy. Software developers have used varying copy protection schemes. These schemes generally prevent a user from copying software or using an unauthorized copy of software. Schemes include \"dongles\", \"code-wheels\", storing the software on difficult to copy storage devices such as copy protected disks, using enabling access codes provided in user manuals or by the developer that must be entered before the software will work, or using hardware specific validation codes that allow a piece of software to operate only on one identified computer. Alternatively, a software developer could use a hardware copy protection scheme where an access code is supplied by a unique or difficult to copy hardware device able to be attached to only one computer at a time. Typically, when licensed software is not given an access code, the standard practice is for the software to refuse to operate. See, for example, U.S. Pat. Nos. 4,562,306 (Method and Apparatus for Protecting Computer Software Using an Active Coded Hardware Device) and U.S. Pat. No. 4,748,561 (Method of Protecting Computer Software).\nCopy protection schemes have not been effective in stemming the tide of piracy, mainly because there are many specialized anti-copy products available to \"unlock\" the copy protection. Further, some copy protection schemes prevent consumers from legitimately making backup copies of the software for use when the authorized copy fails or is destroyed.\nAnother approach to prevent piracy of software is to require the authorized user, upon the first start-up of the software, to input the user's name (which is then stored in the software's code) such that the user's name will appear on the initialization screen on each successive re-use in order to identify the user to whom the software is licensed or sold. This practice is intended to discourage authorized software users from allowing their software to be reproduced, knowing that they can be identified readily from an unauthorized copy of the software. However, as the software producer is unlikely to know when the software is illegally copied and is unlikely to be able to track down unauthorized users, such a scheme does not discourage copying. Further, there is no identification of users who pirate pirated software.\nAn approach to reduce piracy prior to purchase but to allow a potential consumer to test a piece of software is to distribute demonstration software without the full functional features of the software program. This permits the user to sample the user interface and certain features of the software without having to purchase the software product \"sight unseen\". If the user likes the software, the user can contact the developer or distributor for the fully functional software upon registration and payment of a license fee or the purchase price. The problem with such demonstration software is that twice the number of software products have to be produced and distributed. One or more storage mediums containing software must be produced and distributed for demonstration and one or more storage mediums containing different software must be produced and distributed for full function operation. This increases the costs of production and distribution. Further, the user is inconvenienced by twice having to contact or visit the distributor and is delayed until the full function software is obtained.\nOne solution to the problem of having to obtain additional fully functional software is to initially provide to the user the full function program but to have only certain files active for demonstration purposes. Then, if the user wants to enable the full function, the user must contact the developer or distributor for registration and payment of the license fee or purchase price in return for an access code. However, once the access code is revealed, unauthorized copies can be made and the access code can be passed on with the unauthorized copy of the software. Therefore, this is not effective protection. Typically the process also requires the overhead of serializing each individual software copy at a central point of manufacture.\nAs a result of the high overhead costs of the traditional means of manufacture and distribution of software, and the significant level of unauthorized reproduction of software, there is a meed in the software industry for a system that takes advantage of the software consumers' ability to manufacture and distribute copies of software using their own equipment while at the same time controlling access to the software so that the person who designed and coded the software receives fees for its use. In other words, there is a need in the industry for a method of manufacture and distribution that invites rather than restricts consumer participation in the manufacture and distribution process but does not reduce the incentive of the original manufacturer to create the software.\nSoftware is not the only product that suffers from these problems. Other developers and distributors in the information industry, such as developers and distributors of databases and other informational products, also have difficulties controlling access to their product where it can be loaded onto a user's own machine. They have a product that individual users find easy to copy and distribute, but which is costly for the developer to initially design, make and efficiently distribute. Developers and distributors of information products desire to harness the user's ability to copy and distribute their product but only where they can obtain a fee for their initial involvement.\nAs a result, there is a significant and substantial need in the information industry to allow developers to sell their information product to the full potential market efficiently, that encourages and takes advantage of consumer participation in the manufacture and distribution of the information products yet has centralized monitoring and control, that permits consumers to evaluate the product prior to the purchase decision, that effectively combats piracy, and that permits developer profit for funding of support and update services."}
{"text": "Many underdeveloped areas may not have access to broadband services such as satellite location services. Infrastructure expenses may prevent improvement in such areas.\nMany aerial platforms such as airplanes and drones or Unmanned Aerial vehicles (UAVs) have been or will be deployed in various regions around the world for broadband access. Such platforms may be deployed in a more cost effective manner than typical satellite networks.\nThus, there is a need for an efficient way to determine location of and/or track elements using aerial platforms."}
{"text": "This invention is directed to the providing of a transcutaneous nerve stimulator which is designed to be utilized in T.E.N.S. (Transcutaneous Electrical Nerve Stimulation) therapy. T.E.N.S. therapy is based on a non-invasive, non-narcotic concept of pain management which is non-addictive, is not subject to abuse, and does not interact with drugs. T.E.N.S. therapy has already proven to be an effective modality in treating the organic pain problems associated with the following conditions: chronic lumbar and cervical strains or sprains, degenerating disc disease, degenerative arthritic disease, neuropathies, neuralgias, post-lumbar laminectomy syndrome, post-thoracotomy syndrome, bursitis, postphlebitis syndrome, phantom limb syndrome, and tension and migraine headaches.\nEarly attempts to suppress organic pain and other neurophysical effects utilizing electrical stimulation occurred as early as about 2,000 years ago when it was discovered that gout apparently cuold be successfully treated by placing the diseased extremities in a tub of water filled with electric eels. Later, headaches were treated using a similar approach. A detailed, scientific investigation was finally conducted by Professor Galvani of the University of Bologna, which investigation is credited with ultimately leading to the development in the 1800's of electrical equipment for suppression of organic pain.\nThe earliest therapeutic devices utilizing electrical simulation for the most part featured a constant amplitude and rate. Examples of the early art are Benz, U.S. Pat. No. 646,793; Raymond et. al., U.S. Pat. No. 872,148; Tibbals, U.S. Pat. No. 1,059,090; and Call, U.S. Pat. No. 1,908,688. A major problem with electrical stimulation therapy was, and still is, accommodation, whereby the nerve being stimulated in effect accommodates itself over time to the electrical charge, such that the effectiveness of the treatment is diminished.\nIt took scientists a long time to discover, and attempt to address the problem. Nemec, U.S. Pat. No. 2,622,601; DiPerma, U.S. Pat. No. 2,624,342; and Gratzl, U.S. Pat. No. 2,771,554 all disclose electrotherapeutic devices with at least one including means to vary the rate, amplitude or pulse width of the generated electrical pulse. However, merely being able to change either the rate, amplitude or pulse width still resulted in the problem of accommodation occurring, unless an individual manually adjusted the controls prior to the occurrence of accommodation. The process was both labor intensive and inefficient, with respect to the quality of the therapy, since maximum pain relief was not being provided.\nIn 1967, a Dr. Sweet at Massachusetts General Hospital developed the first T.E.N.S. unit. The effectiveness of T.E.N.S. therapy is based on its incorporating two major pain control theories. Under the so-called Gate Control Theory, pain can be inhibited and suppressed by \"closing the gate\" on pain signals as such signals arrive at Central Nervous System centers. This theory postulates that by providing electrical stimulation of a sufficiently high amplitude, the electrical signals race up large myelinated fibers faster than the pain signals travel up smaller myelinated or unmyelinated fibers. The neutral impulses transmitting pain information to the brain thus become interrupted, and since the brain fails to receive the pain signals, no pain is perceived. The other theory incorporated in T.E.N.S. units is the Endorphin Theory, also known as the Endogenous Opiate Theory. This theory postulates that the sustained input of T.E.N.S. signals triggers the release of naturally occurring pain making endorphins and enkephalins (morphine-like substances). These natural substances seemingly block pain signals by a mechanism similar to conventional drug therapy, and inhibit pain information from reaching the brain.\nHowever, T.E.N.S. units, like all electrotherapeutic devices, have suffered problems with accommodation. For example, Geerling, U.S. Pat. No. 4,019,519 issued in 1977, disclosed a unit having only its amplitude adjustable. Miller, U.S. Pat. No. 4,084,595 issued in 1978, disclosed a unit having its rate, amplitude, and pulse width all independently, manually adjustable. However, even this advancement resulted in a less than efficient treatment of the problem of accommodation, since either the therapist or the patient had to, in theory, repeatedly adjust the controls in an attempt to avoid accommodation.\nAlthough variation enables one to deal with the problem of accommodation, pain relief is sacrificed. This is in part due to the interaction between amplitude and pulse width. There is a clinical correlation between amplitude and pulse width with regard to the efficacy of the stimulus. As one shortens the duration of a pulse, its amplitude must be increased to maintain the efficacy of the stimulus. This relationship when plotted graphically is known as a strength-duration curve. Thus not only must the ideal T.E.N.S. units have adjustable amplitude and pulse width, but it must also be able to modulate those values in such a way as to approximate the A-gamma-motor strength-duration curve.\nIn addition to amplitude and pulse width approximating the strenght-duration curve, the rate of the pulses must vary so as to eliminate any potential for accommodation. This explains the short-comings in Reiner, U.S. Pat. No. 2,808,826 which disclosed a unit which permitted instantaneous changes in pulse width and amplitude to two pre-set points along the strength-duration curve, and Maurer, U.S. Pat. No. 4,340,063 which disclosed a unit having its amplitude modulate in response to modulations in pulse width so as to approximate a portion of the strength-duration curve. The rate in Maurer was adjustable, but only to the extent taught by Miller, such that the problem with accommodation still existed."}
{"text": "The present invention is generally directed to a device for injecting heated residual oil into a coke drum as part of a delayed coking process. More particularly, the present invention is directed to a retractable center feed injection device which allows feeding the residual oil into the coke drum from a bottom and center of the coke drum using side-feed entry.\nDelayed coking is one process by which petroleum coke can be produced. As a brief overview, the delayed coking process involves heating byproducts of the oil refining process, known as residual oil, in a furnace to its thermal cracking temperature of approximately 485° to 505° C. This heated residual oil is then pumped into a coke drum where the heat and pressure that build up within the coke drum complete the thermal cracking process. As the heated residual oil cools, petroleum coke is formed which can then be removed from the coke drum.\nIt is typically desirable to feed the heated residual oil from the bottom and center of the coke drum. One reason for this is to minimize the uneven heating of the walls of the coke drum that may otherwise occur if the residual oil is fed towards one side of the wall. To accomplish feeding the residual oil from the center, various so-called “center feed” systems have been developed. These center feed systems employ an injection nozzle which can be inserted into and retracted from the center of the coke drum. Examples of center feed systems are described in U.S. Pat. Nos. 8,702,911 and 8,545,680.\nThese prior art center feed systems employ an injection nozzle that is contained within an inlet sleeve through which the heated residual oil is fed into the injection nozzle. In other words, the combination of the inlet sleeve and the injection nozzle forms the feedline. As a result, all components, including those used to extend and retract the injection nozzle must be configured to withstand the high temperature and pressure of the heated residual oils that they carry and are therefore relatively costly. Additionally, due to the harsh environment in which these components operate, they can require frequent maintenance leading to downtime in the delayed coking process."}
{"text": "Device-to-device (D2D) communications may be utilized for various purposes, such as public safety communications. D2D communications may be associated with standardized technologies, such as LTE, IEEE, etc. In LTE systems, access control and/or priority handling may be used to arbitrate access to and/or usage of wireless resources by terminals."}
{"text": "Field of the Invention\nThe invention relates to an insecticidal water-in-oil (W/O) formulation with at least one insecticidal active substance and at least one burning salt and to the preparation of this formulation. The formulation according to the invention is particularly suitable for the treatment of suitable supports, in particular of paper supports, in an economical one-step process with the aid of conventional application processes. In addition, the present invention relates to insecticidal, smoulderable products which can be prepared by treating a support with the formulation according to the invention.\nDescription of Related Art\nUS-A-2009/0163582 describes a water-in-oil (W/O) formulation in which a polyglycerol fatty acid ester or a sorbitan fatty acid ester is employed as emulsifier, a pesticide, for example a pyrethroid, as active substance and acetyl ester, a methyl ester, an acetyl tributyl citrate, a white mineral oil or a combination of these as the solvent. This formulation takes the form of an aerosol formulation for the end user.\nWO-A-2007/131679 describes a paper impregnated with an insecticidal active substance. The paper here is prepared in a two-step process in which the paper is first pretreated with an aqueous 6% strength potassium nitrate solution and then dried, and then, in a further step, impregnated with active substance solution. WO-A-2007/131679 does not disclose any formulation in which the paper can be treated in a one-step process, in particular with an insecticidal active substance and the potassium nitrate salt.\nWO-A-2011/092722 describes a paper impregnated with an insecticidal active substance, which paper, again, is prepared in a two-step process. Analogously to the process of WO-A-2007/131679, the paper is first treated with a potassium nitrate before being dried and then impregnated with the active substance."}
{"text": "Aspects of the disclosure relate to electrical computers, digital processing systems, and multicomputer data processing. In particular, one or more aspects of the disclosure relate to enabling multicomputer processing and authentication of user data received from a telephone call forwarding system.\nAs computer systems are increasingly utilized to provide automated and electronic services, such computer systems may obtain and maintain increasing amounts of various types of sensitive information, and ensuring the safety, security, and accuracy of this information may be increasingly important. In many instances, however, it may be difficult to ensure the safety and security of such information while also optimizing the efficient and effective technical operations of the computer systems that process such information and/or provide such automated and electronic services."}
{"text": "The present invention is directed to the transmission of data to and from a computer, and more particularly to a system for performing real-time signal processing of data that is serially transmitted to and from a computer.\nVarious devices are known for transmitting data between a computer and a remote site via wide-area telecommunications networks. One of the most widely used devices of this type is the modem, which enables data to be transmitted to and from a computer over a wide-area analog telephone network. Generally speaking, the modem includes one or more sets of registers, typically embodied in an UART or an USART, for storing bits of digital data transmitted to or from the computer, a processor for implementing modem operations, such as dialing a telephone number or answering a ringing signal, in response to commands sent from the computer and stored in the UART, and a modulator/demodulator for converting digital bits of data to be transmitted into analog signals, and vice versa. Originally, all of these features were hardwired in a separate peripheral device that could be externally connected to the computer via a serial I/O port, or internally connected to the computer\"\"s data bus. More recently, some of the functions associated with these features, most notably the processing of commands to implement modem operations, have been removed from the hardwired configuration and incorporated into the computer itself. This approach has increased the versatility of the modem. For example, while the hardwired modem configuration had to be specifically designed for the telephone system requirements of a particular country, the later approach could enable a single product to be used in a variety of countries, each of which might have different telephone signaling requirements. Similarly, since the computer itself was handling the data to be transmitted, additional services, such as the ability to send information as a facsimile transmission, in which graphical data is processed, became feasible. However, the presence of the UART, or similar such device through which the data must flow, still limits the effective rate at which the data can be exchanged between the computer and the telephone lines.\nTo enhance the performance of modems, a digital signal processor (DSP) has been incorporated into its structure. In this arrangement, the modem software was designed to cooperate with the DSP to provide data thereto for processing prior to transmission or after reception over the telephone line. While the addition of the DSP provided increased capabilities in terms of the speed at which the data could be transmitted over a telephone network and the ease with which the modem could be configured, it was still limited in the types of data that could be processed. More particularly, because of the restrictions imposed by passing the data through an UART or the like, even the most modem modems are only capable of effectively transmitting data in the range of 9.6-14.4 Kb/sec. While this rate of data transfer may be useful for transmitting static information such as text files or the like, it is not suitable for many real-time applications in which the data is provided at much higher rates, such as speech or video conferencing.\nFurther in this regard, the modem control software had to be designed to work with the specific DSP incorporated into the computer. If a different DSP was to be used, the modem control software had to be reprogrammed to work with the new DSP.\nWhile the analog telephone network was the only practical medium for transmitting information between geographically distributed computers for quite some time, more recently other, non-analog transmission mediums have become available. Examples of these other mediums include the integrated services digital network (ISDN), private branch exchange (PBX) telephone systems, and TI digital data links. Since information is transmitted over these mediums as digital data, conventional analog modem circuits are not suited for use with them. Thus, for example, a standard Group III facsimile machine cannot operate on a digital PBX system.\nSimilarly, digital signal processing systems which are designed to work with PCM-encoded analog data that is received and transmitted via a modem are likewise not suited for use with these other types of transmission mediums. While it is possible to incorporate another DSP system into a computer that can handle data transmitted via any of these digital networks, it would be more desirable to provide a single system that can process data that is received over any type of transmission medium, whether it be digital or analog. Further in this regard, it is desirable to provide a signal processing system that is not limited to one specific DSP, but rather one that can operate with any of a variety of different types of signal processors.\nWhen personal computers communicate with one another through non-modem transport facilities, they typically operate in a burst mode. While this mode of operation enables data to be transferred at much higher rates than with modems, it is still not suitable for applications such as video or speech processing. These types of applications require isochronous data handling, i.e. data that is transmitted at a constant bit rate and that must be processed in real time. Generally speaking, therefore, it is desirable to provide a serial data transmitting and receiving system that is capable of processing real-time isochronous data.\nFurther in this regard, it is desirable to provide such a system that is capable of handling streams of data pertaining to different functions. For example, in a video conferencing application, speech data is transmitted at the same time as video and other graphic information. However, each of these types of data must be processed separately in real time. It is desirable, therefore, to provide a data transfer system that can handle each of two or more types of data at isochronous rates.\nThe present invention provides a data transmission system having a real-time data engine for processing isochronous streams of data. An interface device provides a physical and logical connection of the computer to any one or more of a variety of different types of data networks, including analog telephone, ISDN, PBX and the like. Data received at this device is presented to a serial driver, which disassembles different streams of data for presentation to appropriate data managers, such as the operating system of the host computer, a service provider or an application program. A device handler associated with the interface device sets up data paths and issues service requests. The device handler also presents data and commands from the data managers to a real-time data processing engine, that can be used for a variety of applications such as voice recognition, speech compression, and fax/data modems. This real-time engine can be shared by any application program running on the host computer.\nThe invention enables any arbitrary type of data, such as voice, facsimile, multimedia and the like, which is transmitted over any type of communication network, to be handled with any type of real-time engine, by abstracting the functions of each of the elements of the system from one another. This abstraction is provided through suitable interfaces that isolate the transmission medium, the data managers and the real-time engine from one another. The data is provided to the real-time engine as multiple streams of isochronous data, i.e. it is delivered as it arrives over the network without data handling delays. This feature allows more complex applications, such as speakerphones, videophones and high-speed modems to be readily implemented.\nThese features of the invention, as well as the advantages offered thereby, are described in greater detail hereinafter with reference to a specific embodiment illustrated in the accompanying figures."}
{"text": "A polybutylene terephthalate resin is excellent in an injection molding property, mechanical characteristics, a heat resistance, electric characteristics, a resistance to chemicals, etc. Utilizing these properties, the polybutylene terephthalate resin has been widely utilized as products by an injection molding in fields of automobile parts, machine parts, electro-communication parts, etc. Further, it is also excellent in a mold transfer property and is advantageously used also in lamp members such as extension uses of automobiles for which a good external appearance is particularly demanded. Furthermore, in view of its use, it is necessary to highly control the heat resistance of the resin, and a suppression of gas generation (low gas emission) upon the molding, etc.\nOn the other hand, an LED light has been installed for headlamps of high-class automobiles in recent years and a lamp design has been also started in a drastic change from the former one. For example, in a reflector type (a type wherein a light from a light source is irradiated by means of reflection using a reflector), it is necessary that an area around the light source is metallized with aluminum. Recently, a new design has been appearing wherein the reflector type is changed to a projector type (a type wherein the light from the light source is irradiated by condensing to a front lens) and the area around the light source is changed to a black colored design. However, due to this change, there has been resulted a problem of condensation of a sunlight. That is, when the sunlight is reflected by a projector lens and is condensed to a surrounding black area, a temperature of the condensed area becomes very high whereby a surface melting occurs. (Such a phenomenon itself has been known from old time already but, up to now, it has not caused any particular problem in view of a lamp design.) Taking such a problem into consideration, there has been a demand for materials which transmit infrared light so that the temperature does not increase even by the condensation of the sunlight and also for materials which have a heat resistance resulting in no surface melting even at the high temperature.\nAs to a technique for transmitting the infrared light, there are disclosures, for example, in Patent Documents 1 to 3 for a resin composition containing a resin consisting of polybutylene terephthalate or of polybutylene terephthalate and a polybutylene terephthalate copolymer together with a noncrystalline resin such as a polycarbonate resin, a styrene acrylonitrile resin, a polyester resin containing 1,4-cyclohexanedimethanol ingredient, etc. However, although this technique is useful for enhancing the infrared light transmittance of the resin of a polybutylene terephthalate type, a heat deflection temperature significantly lowers by addition of the noncrystalline resin. Accordingly, it is difficult to practically use this resin composition particularly as a lamp part.\nIn addition, as to a technique for preparing a black colored polyester resin composition which suppresses a temperature rise due to irradiation of a sunlight, there is disclosed, for example, in Patent Document 4 for a polyethylene terephthalate resin composition which does not contain carbon black but contains a pigment which is made into black color by mixing several pigments. According to this invention, it is possible to lower a temperature corresponding to a heat quantity accumulated by an infrared light-absorbability of carbon black, as a result of the use of the pigment of a non-carbon black type. However, an effect of suppressing the temperature rise thereby is small and, accordingly, there is yet a room for improvement."}
{"text": "To satisfy the increasing demand of bandwidth and speed, pluggable transceiver modules are used in line cards on various network devices (e.g., switches, routers, etc.). The pluggable transceiver modules are used to convert electrical signals to optical signals or in general as the interface to a network element copper wire or optical fiber. There are tests that need to be performed on the interface between the transceiver and a host for compliance with standards and to ensure that the transceiver will work properly with the host. Conventional testing is time intensive and involves high costs, and may not be scalable in a manufacturing and production environment.\nCorresponding reference characters indicate corresponding parts throughout the several views of the drawings."}
{"text": "The vascular endothelial growth factors (VEGFs) are a major family of angiogenic proteins involved in endothelial cell activation, proliferation, and survival, particularly during retinal proliferative diseases and tumorigenesis. VEGFs belong to the VEGF-PDGF (platelet-derived growth factor) super-gene family, and are small glycoprotein dimers that bind receptors expressed on vascular and lymphatic endothelial cells. There are currently seven known ligands in the VEGF family: VEGF-A (VEGF), VEGF-B, VEGF-C, VEGF-D, VEGF-E (viral-derived), and placental growth factor (PlGF)-1 and -2. These VEGF ligands mediate their effects by binding to one or more of the three known VEGF receptors (VEGF-Rs), each of which possess receptor tyrosine kinase activity. VEGF-R1 (Flt-1) is predominantly expressed on endothelial cells and monocytes, binds VEGF and VEGF-B, and appears to mediate endothelial and monocyte migration. VEGF-R2 (i.e., human KDR or murine Flk-1) is mainly expressed on endothelial cells, is selective for VEGF (and particular fragments of VEGF-C and VEGF-D), and mediates VEGF-induced endothelial cell proliferation, survival, and migration, as well as vascular permeability. VEGF-R3 (Flt-4) is mainly expressed on lymphatic endothelial cells and binds VEGF-C and VEGF-D to promote lymphangiogenesis. VEGF-R1, -R2, and -R3 are each expressed on some tumor cells. Binding of VEGF to the VEGF receptors triggers receptor dimerization, leading to subsequent receptor activation and signal transduction. VEGF binding to VEGF-R2 initiates a signal transduction pathway that is dominant in promoting angiogenesis. This pathway involves receptor activation with subsequent induction of intracellular signaling. Receptor activation in this case entails three basic events: (i) VEGF binding to VEGF-R2, (ii) receptor dimerization, and (iii) receptor autophosphorylation (and hence activation) of the receptor tyrosine kinase. Intracellular messengers such as phospholipase C and phosphatidylinositol-3-kinase bind directly to the autophosphorylated form of VEGF-R2 and become phosphorylated by the receptor tyrosine kinase, which subsequently triggers an intracellular cascade of signaling events leading to nuclear signals that ultimately promote cell proliferation, migration, and survival (anti-apoptosis), and increase vascular permeability.\nAberrant angiogenesis is associated with a variety of disease states, including cancer (Holash 2002). VEGF signaling has been verified to play a role in both normal vascular development and the pathologic angiogenesis associated with various diseases (Erikkson 1999; Ferrara 1999; Yancopoulos 2000). VEGF promotes vascular endothelial cell growth and increases vascular permeability (Ferrara 2004).\nPrevious studies have revealed elevated VEGF expression levels in a majority of tumor types (Berkman 1993; Brown 1993; Brown 1995; Dvorak 1995; Mattern 1996). Studies have also revealed increased VEGF levels in subjects with ocular angiogenic diseases such as wet AMD. Wet AMD accounts for only around 10% of total AMD cases, but causes approximately 90% of blindness arising from AMD. Wet AMD is characterized by choroidal neovascularization (CNV), the development of abnormal blood vessels beneath the retinal pigment epithelium layer of the retina. VEGF-A is believed to play a major role in the formation of these vessels, which leak beneath the macula and cause retinal distortion and vision deterioration.\nTherapeutic anti-VEGF antibodies are currently available. For example, the humanized IgG1 monoclonal antibody Bevacizumab (a.k.a., Avastin®, sold by Genentech, San Francisco, Calif.; also referred to herein as BM-1) binds human VEGF with an affinity (KD) of approximately 500 pM. While Bevacizumab has been used to treat a variety of cancers, there is a need in the art for antibodies with greater in vivo efficacy. Such antibodies present significant technical challenges and are highly elusive because merely increasing the binding affinity of a VEGF antibody does not necessarily increase its in vivo efficacy (Liang 2006)."}
{"text": "This invention relates to the creation of Bragg gratings in photosensitive optical waveguides. It is known that such a grating can be created by illuminating the photosensitive waveguide from the side, writing the lines simultaneously with an interferometrically generated grating fringe pattern of light. Such a fringe pattern can be created using two-beam interferometry, or as a fringe pattern generated in the vicinity of a diffraction grating, typically a phase grating, through which light is caused to pass.\nTwo broad regimes for producing gratings in photosensitive optical waveguides are known, respectively characterised as type I and type II. A type I grating is created by relatively low fluence exposure to a fringe pattern, and the refractive index change that is produced grows towards a saturation level with increasing fluence, which is normally, but not necessarily, provided on a cw basis. The manufacture of such type I gratings is for instance described in U.S. Pat. No. 4,474,427 (K O Hill et al) and U.S. Pat. No. 4,725,110 (W H Glenn et al). A type II grating is created by significantly higher fluence exposure to a fringe pattern, and the effect has been attributed to some `damage` effect upon the core of the fibre. The manufacture of such a type II grating is for instance described by J L Archambault in a paper entitled `High Reflectivity and Narrow Bandwidth Fibre Gratings Written by Single Excimer Laser Pulse`, Electronics Letters (7 Jan. 1993), Vol. 29 No. 1, pp 28-29. Published papers on type II grating creation have recited the use of ultra-violet radiation at 248 nm from a KrF excimer laser, the fluence being confined to a single pulse typically of about 20 ns duration. The suggestion that the observed effect is a `damage` effect is based upon observation of the effects of such fluences incident upon fibre preforms similar to those from which the optical fibres themselves were drawn. Damage is inferred when ablation from the surface of preforms is detected using a sensor that directs an interrogation beam of laser light laterally through a zone just above the point where the 248 nm ultra-violet light is incident upon the preform surface. When using a KrF excimer laser in this way, it is found that there is a relatively narrow safety margin between the fluence necessary to reach the lower threshold of the onset of refractive index modifying `damage`, and the fluence necessary to reach the upper threshold at which catastrophic disruption of the fibre is liable to occur.\nFor the purpose of this invention the term `damage`, as used in the context of Bragg grating creation, is defined to cover a process of creating a Bragg grating in a photosensitive waveguide using a fluence sufficient to produce ablation effects when applied to an optical fibre preform having the same optical core composition as that of the waveguide."}
{"text": "The present invention relates to orthopedic treatment processes and, in particular, the present invention relates to preparing orthopedic treatment protocols used in conjunction with computerized or digitalized orthopedic treatment devices.\nOrthopedic treatment historically involved a physician examining and diagnosing an orthopedic injury in a patient, prescribing a treatment protocol of activities or exercises for the patient to follow in order to facilitate healing, and subsequent re-examination to assess patient progress. Additionally, the patient was traditionally guided and assisted in following the prescribed treatment protocol by trained medical professionals, such as physical therapists, who could inform and advise the attending physician concerning patient compliance with the protocol and communicate and assist with the patient to provide desired activity details and elicit patient response. The traditional treatment path often included either hospitalization or patient visits at a physical therapy facility.\nIn modern times, financial pressure upon the medical arts and the surrounding medical industry has increased the number of patients each physician must treat and reduced the rate of hospitalization and tendency to employ physical therapy facilities, as well as reduced the direct supervision of the patient activities by the physical therapist. Computerized devices have been developed that at least partially substitute for the physical therapist contact, and yet monitor patient activities under a treatment protocol. One particularly innovative device system, the IZEX sensor-instrumented orthosis and associated Smart IDEA(trademark) computer/conununicator, not only replace much of the physical therapist\"\"s function of (1) advising and instructing the patient and (2) advising the attending physician of patient compliance, but also allow an improved measuring and monitoring of patient rehabilitation activities and exercise parameters, such as effort exerted in rehabilitation exercises or stress applied to the orthopedic injury. This improved monitoring enables exploitation of a long observed and literature documented phenomenon of improved recovery in response to appropriately applied exercises to orthopedic injuries. The topic of accelerated and improved recovery through the use of controlled bio-feedback based rehabilitation has been reviewed extensively by one of the present inventors in patents U.S. Pat. No. 5,052,375; U.S. Pat. No. 5,368,546; U.S. Pat. No. 5,484,389; and U.S. Pat. No. 5,823,975, and the entire disclosures of these patents are incorporated herein by reference.\nThe ultimate goal of efficiency and optimal and accelerated recovery outcomes still remains elusive, however, since the utilization of the IZEX(trademark) orthosis brace system and Smart IDEA(trademark) computer/communicator continues to rely upon physician examination, diagnosis and prescription of a treatment protocol for the injured patient. The physician may not readily know nor have available information concerning the optimal treatment protocol for an accurately diagnosed injury. It would be a significant advance in orthopedic treatment if a physician or other treatment professional could be rapidly advised concerning optimal treatment information based upon up-to-date experiential outcomes of similar treated injuries. It would also be a significant advance if the physician or treatment professional could leverage their own expertise and their colleagues\"\" most recent knowledge to appropriately modify and adapt previously successful protocols to fit a new patient. It would also be a significant advance if the protocol could be installed in a handheld computer (monitoring device/computer/communicator) device with ease and efficiency. Additionally, it would be a significant advance to allow modification, particularly real-time modification, of rehabilitation exercise protocols by a user or in response to a user request, with such modification being limited by reasonable constraints. The following invention provides such advances to the orthopedic arts."}
{"text": "As compared with red sulfide low voltage electroluminescent phosphor Zn1-xCdxS (x=0˜1.0), titanate substrate has a good chemical stability, and the phosphor obtained by doping rare earth luminescent center ions, such as CaTiO3, has a better color purity. The coordinates of red color of Pr3+ excited by photoluminescence and cathode ray are: x=0.680, y=0.311, which are very close to that of the ideal red according to NTSC color gamut. Considering the material stability and luminous colors, the titanate substrate phosphor activated by rare earth ion is expected to replace red sulfide phosphor Zn1-xCdxS and become a new generation of non-toxic, highly stable red FED phosphor material. Accordingly, it is the goal of researchers for further improving the luminescent properties of this material."}
{"text": "Soft tissue replacement methods are commonly used for a wide variety of clinical and cosmetic purposes. One use involves reconstructive applications that rebuild and restore a body part or structure to correct deformities from congenital anomalies, trauma, cancer, infections, disease, or medication side effects. The replacement tissue serves to support surrounding tissue and to maintain the normal appearance of the body. The restoration of this normal appearance has an extremely beneficial psychological effect on post-operative patients, alleviating much of the shock and depression that often follows extensive surgical procedures. Another use involves augmentation applications that alter a body part or structure usually to improve its cosmetic or aesthetic appearance. Augmentation of the appearance also has beneficial psychological effects that improve self-esteem, well-being, and confidence of an individual. A third use involves structural applications that provide support to a body part or structure to improve its function, thereby alleviating a symptom associated with a disorder involving soft tissue loss. Examples of such disorders include, without limitation, stress urinary incontinence, fecal incontinence, vocal cord paralysis, vocal atrophy, vocal implantation, intubation trauma, post-hemilaryngectomy defects, irradiation damage, lumbar disc repair, and plantar footpad repair.\nSoft tissue replacement methods currently rely on two general approaches: 1) implantation of artificial or alloplastic fillers like soft tissue implants, and injectable polymers and hydrogels; and 2) transplantation of tissue like tissue flaps and autologous tissue transfers. A drawback to the use of artificial or alloplastic fillers it that these inorganic materials lack any metabolic activity, do not become physiologically incorporated into the body, and as such, surrounding tissue and blood supply does not develop within the implanted material. In addition, artificial or alloplastic fillers risk migration and/or extrusion from the implant site. Furthermore, many of these fillers produce only temporary effects because the body rapidly reabsorbs them. Fillers providing long-term effects frequently induce a foreign body response resulting in formation of an avascular, fibrous capsule around the filler, which limits performance, distorts the aesthetic appearance of the surrounding area, and can cause pain to the individual. Furthermore, long-term fillers, like soft tissue implants, do not remodel with the aging tissue resulting in a material that may not be aesthetically acceptable as the patient's tissues undergo the normal physiologic changes associated with aging.\nAlthough the use of transplanted tissue avoids the problems associated with artificial or alloplastic fillers, drawbacks are also associated with these procedures. In these procedures, loss of transplanted tissue volume over time as a result of its resorption by the body is a major problem. For example, transplantation of adipose tissue generally results in a loss of 20% to 90% of its volume within the first year. This tissue loss is unpredictable and is a result of poor survival and/or regeneration from progenitor cells in the transplanted tissue due to necrosis and a lack of vascular formation. With respect to adipose tissue, tissue breakdown is associated with traumatic rupture of the cells, avascular necrosis, apoptosis of the adipocytes, inflammation secondary to apoptosis, fibrosis and contraction of the graft, and/or delipidation of the adipocytes with subsequent volume loss. Failed tissue grafts sometime produce stellate and irregular nodules with calcifications. As such, transplanted tissue methods are usually performed two or three times to obtain the desired effect, resulting in massive time and cost.\nOne of the major underlying causes for tissue breakdown seen in transplanted tissue methods is the lack of a blood supply sufficient to support the transplanted tissue. For example, alleviation of tissue ischemia is critically dependent upon formation of new blood vessels. The growth of new blood vessels and associated vasculature or repair or remodeling of existing blood vessels and associated vasculature provides collateral circulation in and around an ischemic area, improves blood flow, and alleviates the symptoms caused by the ischemia. Thus, compositions and methods that promote new blood vessel formation within a transplanted tissue are needed as this will improve the survival rate of such tissues, thereby achieving reliable long-term survival of grafted tissues.\nThe present specification provides novel compositions and soft tissue replacement methods using compositions that reduce tissue volume loss by increasing the survival rate of the transplanted tissue. This improved survival rate is achieved by administering a compound that promotes new blood vessel formation, thereby ensuring that a blood supply adequate to support the transplanted tissue is established."}
{"text": "The present invention relates generally to a musical pickup arrangement for the conversion of sound vibrations into electrical oscillations, and more specifically, to a musical pickup arrangement incorporating a condenser-type microphone for converting sound vibrations generated in a musical instrument into electrical signals that can be supplied to a loud-speaker.\nAs is known, microphones and other transducers have been employed in connection with musical instruments to convert sound vibrations into electrical oscillations. By means of suitable electronic apparatus, these electrical oscillations have been amplified and reconverted into audible sounds. In certain cases, piezoelectric translating devices have been employed in pianos and stringed instruments for converting mechanical vibrations of the resonator or sounding board of the instrument into electrical oscillations which, in turn, have been amplified and reconverted into audible sounds. These devices have been placed, for example, directly on the sounding board of the musical instrument and at various other locations above the sounding board, as evidenced by the prior art. However, the sound produced by each of these arrangements has generally been poor. The reasons why the sound produced by such transducers is often poor is perhaps due to the fact that the device is secured directly to the sounding board. This may, for example, be due to the fact that the natural frequency of the sounding board is altered when contacted by the pickup.\nIn other cases, contact pickup devices used on musical instruments have been of the magnetic-induction type. In these devices, a vibrating or resonating cavity is in contact with a coiled magnet which transduces the vibrations of the cavity into small electrical impulses, thus creating the transformation of magnetic to electrical energy.\nThe magnetic-induction type of transducer has been quite reliable for many years; however a falsification of the actual color or \"timbre\" of certain instruments is lost when this type of pickup is used on a nylon-stringed guitar or the like. Of the six strings on the classical guitar, the top three are made of pure nylon, similar to a nylon fishing line, but milled to various thicknesses. The bottom three base strings have a nickel alloy wrapping around several fine threadlike nylon strands. A magnetic pickup device used on this type of instrument transforms sound from the bottom three base strings evenly due to their outer conductive coating, but the non-conductive nature of the top three treble strings prevents an even transfer of energy and usually results in sound which is uncharacteristic, or \"tinny\"."}
{"text": "Communication systems are known to support wireless and wired communications between wireless and/or wired communication devices. Such communication systems range from national and/or international cellular telephone systems to the Internet to point-to-point in-home wireless networks. Each type of communication system is designed, and hence operates, in accordance with one or more communication standards. For instance, wireless communication systems may operate in accordance with one or more standards including, but not limited to, IEEE 802.11, Bluetooth, advanced mobile phone services (AMPS), digital AMPS, global system for mobile communications (GSM), code division multiple access (CDMA), local multi-point distribution systems (LMDS), multi-channel-multi-point distribution systems (MMDS), and/or variations thereof.\nDepending on the type of wireless communication system, a wireless communication device, such as a cellular telephone, two-way radio, personal digital assistant (PDA), personal computer (PC), laptop computer, and/or home entertainment equipment, communicates directly or indirectly with other wireless communication devices. For direct communications, also known as point-to-point communications, participating wireless communication devices may tune their receivers and transmitters to the same channel, or channels, such as one or more of the plurality of radio frequency (RF) carriers of the wireless communication system, and communicate over one or more of those channels. For indirect wireless communications, each wireless communication device communicates directly with an associated base station for use of cellular services, for example, and/or an associated access point for use of an in-home or in-building wireless network, for example, via an assigned channel, or channels. To complete a communication connection between the wireless communication devices, the associated base stations and/or associated access points communicate with each other directly, via a system controller, via the public switch telephone network, via the internet, and/or via some other wide area network.\nFor each wireless communication device to participate in a wireless communications session, it may include a built-in radio transceiver comprising a receiver and transmitter, or it may be coupled to an associated radio transceiver, such as a station for in-home and/or in-building wireless communication networks, and/or an RF modem. The transmitter may be adapted to convert data into RF signals by modulating the data in accordance with the particular wireless communication standard to produce baseband signals and mixes the baseband signal with a local oscillator signal in one or more intermediate frequency stages to produce the RF signals. The radio receiver may comprise an antenna section, a filtering section, a low noise amplifier, an intermediate frequency (IF) stage, an equalization stage, and/or a demodulator. The antenna section may receive RF signals and may provide them to the filtering section, which, in turn, may pass RF signals of interest to the low noise amplifier (LNA).\nThe LNA may amplify the received RF signals of interest and may provide them as amplified RF signals to the IF stage. The IF stage may step down the frequency of the RF signals of interest to an intermediate frequency or to baseband. The IF stage may then provide the intermediate frequency signals, or baseband signals, to the equalization stage. The equalization stage may add loss or delay to specific frequencies to produce a flat frequency response. The equalization stage may then output an in-phase (I) and/or a quadrature (Q) component to the demodulator. The demodulator may be adapted to synchronize the signal and/or to recapture the data in accordance with a demodulation protocol. During operation, the signal quality of a received signal may be utilized to determine how best to receive and demodulate the signal.\nIn some conventional systems, signal quality of an RF signal is usually estimated by using received signal strength indicator (RSSI) measurement and packet error rate (PER). However, due to interference environment, the received signal strength indicator measurement and packet error rate measurement may not provide an accurate estimation of the signal quality. During reception, due to interference, it may be difficult to determine whether a received signal is an intended received signal or interference. For example, a strong signal may be received with good signal quality, but this signal may be interference rather then the intended received signal. Furthermore, to get an accurate estimate of the packet error rate, a large enough population of packets is required and this may take a substantial amount of time.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "1. Field of the Invention\nThis invention relates to a magnetic toner for rendering latent images visible in image-forming processes such as electrophotography, electrostatic recording, magnetic recording and toner jetting, and also relates to an image-forming method making use of the magnetic toner.\n2. Related Background Art\nA number of methods are conventionally known as electrophotography. In general, copies or prints are obtained by forming an electrostatic latent image on an electrostatic latent image bearing member (hereinafter also xe2x80x9cphotosensitive memberxe2x80x9d) by utilizing a photoconductive material and by various means, subsequently developing the electrostatic latent image by the use of a toner to form a toner image as a visible image, transferring the toner image to a transfer medium such as paper as occasion calls, and then fixing the toner image to a recording medium by the action of heat, pressure, or heat-and-pressure.\nApparatus for such image formation include copying machines and printers. In recent years, as printers, LED printers or LBP printers are prevailing in the recent market. As a trend of techniques, there is a tendency toward higher resolution. More specifically, those which hitherto have a resolution of 240 dpi or 300 dpi are being replaced by those having a resolution of 600 dpi, 800 dpi or 1,200 dpi. Accordingly, with such a trend, developing systems are now required to achieve a high minuteness. Copying machines have also made progress to have high functions, and hence they trend toward digital systems. In this trend, chiefly employed is a method in which electrostatic latent images are formed by using a laser. Hence, the copying machines also have come to have a high resolution. Also, with an improvement in image quality, it is much sought to achieve a higher speed and a longer service life.\nIn developing systems used in such printers and copying machines, toner images formed on the photosensitive member in the step of development are transferred to a recording medium in the step of transfer via, or not via, an intermediate member. Any transfer residual toner and fogging toner at non-image areas, left on the photosensitive member are removed in the step of cleaning, and is stored in a waste toner container. In this cleaning step, blade cleaning, fur brush cleaning, roller cleaning and so forth are conventionally used. When viewed from the standpoint of apparatus, the whole apparatus must be made larger in order to provide such a cleaning means. This has been a bottleneck in attempts to make apparatus compact. In addition, from the viewpoint of ecology, a system that may produce no waste toner is long-awaited in the sense of effective utilization of toner. Thus, it is sought to provide a toner having a high transfer efficiency and less causing fog.\nFrom the viewpoint of making apparatus compact, one-component developing systems are preferable because they require no carrier particles such as ferrite particles or iron powder which are required in two-component developing systems. Also, since in the two-component developing systems the concentration of toner in two-component developers must be kept constant, a device for detecting toner concentration so as to supply the toner in the desired quantity is required, resulting in a large size for the developing assemblies. In the one-component developing system, on the other hand, such a device is not required, and hence the developing assemblies can also be made small and light-weight as being preferable. Magnetic toners used in such image-forming processes are commonly chiefly composed of a binder resin and a magnetic material and besides contain additives such as a charge control agent and a release agent which are used to bring out properties necessary as toners. As a colorant of the magnetic toner, the magnetic material is used as it is as the colorant, or a non-magnetic inorganic compound, organic pigment or dye is used together with the magnetic material. As the release agent, used are waxes sparingly compatible with the binder resin, as exemplified by low-molecular weight polyethylene and low-molecular weight polypropylene.\nHowever, developing systems making use of an insulating magnetic toner have a problem concerning the insulating magnetic toner used. The problem is that, in insulating magnetic toner particles, a finely powdery magnetic material is mixed and dispersed in a considerable quantity, and it affects fluidity and triboelectric chargeability of the magnetic toner because magnetic fine particles constituting the magnetic material stand partly uncovered to the surfaces of toner particles, consequently causing variation or deterioration of various performances required for the magnetic toner, in relation to developing performance and running performance of the magnetic toner. This is presumed to be due to the fact that magnetic fine particles having a relatively lower electrical resistance than the resin constituting the magnetic toner particles are present at the surfaces of the magnetic fine particles. Also, the chargeability of the magnetic toner also has a great influence on development and transfer, and is closely concerned with image quality. Accordingly, it is sought to provide a magnetic toner which can stably provide a high charge quantity.\nTo cope with this problem, proposals concerning magnetic iron oxides to be contained in magnetic toners are hitherto made, but there is room for further improvement.\nFor example, Japanese Patent Application Laid-Open No. 62-279352 discloses a magnetic toner containing a magnetic iron oxide incorporated with silicon element. In such a magnetic iron oxide, the silicon element is intentionally brought into existence inside the magnetic iron oxide, but there is room for further improvement in the fluidity of the magnetic toner containing the magnetic iron oxide. Japanese Patent Publication No. 3-9045 discloses adding a silicate to control the shape of magnetic iron oxide to be spherical. In the magnetic iron oxide thereby obtained, the silicon element is rich distributed inside the magnetic iron oxide fine particles because of the use of the silicate for the controlling of particle shape of magnetic fine particles and the silicon element is less present at the surfaces of the magnetic iron oxide fine particles, thus, because of a high smoothness of the magnetic iron oxide fine particles, the fluidity of the magnetic toner can be improved to a certain extent. However, it is preferable to more improve the close adhesion between the binder resin constituting magnetic toner particles and the magnetic iron oxide. Japanese Patent Application Laid-Open No. 61-34070 discloses a process for producing triiron tetraoxide by adding a hydroxosilicate solution to triiron tetraoxide in the course of oxidation reaction. The triiron tetraoxide fine particle obtained by this process has silicon element in the vicinity of its surface, but the silicon element is present in layer in the vicinity of the surface of the triiron tetraoxide fine particles. Hence, there is a problem that the surface is weak to mechanical shock such as friction.\nMeanwhile, toners are produced by melt-mixing a binder resin, colorant and so forth and uniformly dispersing them, followed by pulverization by means of a fine griding mill and then classification by means of a classifier to obtain toners having the desired particle diameters (pulverization process). To make toners have fine particle diameters, there is a limit to the range of material selection. For example, colorant-dispersed resin compositions must be brittle enough to be pulverizable by means of economically available production apparatus. Since the colorant-dispersed resin compositions are made brittle because of such a requirement, particles having particle diameters in a broad range tend to be formed when such compositions are actually pulverized at a high speed, so that, in particular, fine particles (particles having been pulverized in excess) having a relatively large proportion are formed in a large quantity and also the magnetic fine particles tend to come off from the resin during pulverization. Moreover, such highly brittle materials tend to be further pulverized or powdered when used as developing toners in copying machines or printers.\nAs a countermeasure therefor, Japanese Patent Application Laid-Open No. 2-256064 disclose, in the production of pulverization toners, a magnetic toner production process in which magnetic fine particles standing free are removed by classification after pulverization. In pulverization processes, however, the magnetic iron oxide fine particles essentially come to stand uncovered to the surfaces of magnetic toner particles, and hence a problem tends to arise in fluidity of magnetic toner particles and in charging stability in severe environment, resulting in a low transfer performance. Thus, there is further room for improvement.\nIn the pulverization process, it is also difficult to uniformly disperse solid fine particles such as a magnetic powder and a colorant in the resin. Depending on the degree of such dispersion, this can be one of the causes of an increase in fog and a decrease in image density.\nIn the pulverization process, making toner particles finer in order to achieve high minuteness and high image quality accompanies a lowering of uniform chargeability and fluidity of toners.\nIn order to overcome such problems of toners which are ascribable to the pulverization process and further to satisfy requirements stated above, proposed are processes for producing toner particles by suspension polymerization.\nToner particles produced by suspension polymerization (hereinafter xe2x80x9csynthetic toner particlesxe2x80x9d or xe2x80x9csynthetic tonerxe2x80x9d) can readily be produced in fine particles. In addition, toner particles obtained have a spherical shape and hence have a superior fluidity. This is advantageous for achieving a high image quality.\nHowever, incorporation of magnetic fine particles into such synthetic toner particles tends to make the toner particles have a low fluidity and a low charging performance, tending to cause a lowering of developing performance. This is because the magnetic fine particles are commonly hydrophilic and hence tend to be present at toner particle surfaces in the suspension polymerization, which makes use of an aqueous medium. This is also because, in the step of granulation carried out when magnetic synthetic toner particles are produced, hydrophilic magnetic fine particles may partly move to the aqueous medium to become present as free magnetic fine particles having come off from the magnetic toner particles. In order to solve this problem, it is important to modify the surface properties the magnetic fine particles have.\nTo improve dispersibility and enclosure property of the magnetic fine particles in the synthetic toner particles, proposals are made in a large number in regard to surface modification of the magnetic fine particles. For example, Japanese Patent Applications Laid-Open No. 59-200254, No. 59-200256, No. 59-200257 and No. 59-224102 disclose techniques for treating magnetic fine particles with silane coupling agents of various types, and Japanese Patent Application Laid-Open No. 63-250660 discloses a technique for treating silicon-element-containing magnetic fine particles with a silane coupling agent.\nSuch treatment brings about a certain improvement in the dispersibility in the magnetic toner particles. However, the magnetic fine particle surfaces must be made uniformly hydrophobic, and the uncovering of the magnetic fine particles to the magnetic toner particle surfaces must be more controlled.\nMeanwhile, with regard to the quantity of magnetic fine particles at the magnetic toner particle surfaces, a toner having a special structure in which any magnetic fine particles are not present in toner particle surface layers is proposed as disclosed in Japanese Patent Application Laid-Open No. 7-209904. This toner is advantageous in that it ensures superior enclosure of magnetic fine particles and can be free from any uncovering of the magnetic fine particles to the magnetic toner particle surfaces. However, such a toner must be produced by a complicate process, and can be produced with difficulty in an industrial manufacture scale. Also, its repeated use over a long period of time in an environment of low humidity may cause a lowering of image quality which is due to charge-up of the magnetic toner. Thus, more improvement has been necessary for the stability of charging of the magnetic toner.\nFurther, a technique for making the particle diameter of the magnetic toner particles smaller to achieve higher image quality is disclosed in Japanese Patent Application Laid-Open No. 1-112253. However, with the magnetic toner particles having such a smaller diameter, it is more difficult to obtain the uniform dispersion of the magnetic powder and the enclosure thereof, tending to give rise to the above-mentioned various problems.\nFor the purpose of improving the fluidity and charging performance of toners, methods are also proposed in which inorganic fine powder is added as an external additive, and are put into wide use. For example, Japanese Patent Applications Laid-Open No. 5-66608, No. 4-9860 and so forth disclose external addition of an inorganic fine powder having been subjected to hydrophobic treatment or an inorganic fine powder having been subjected to hydrophobic treatment and thereafter further to treatment with a silicone oil. Japanese Patent Applications Laid-Open No. 61-249059, No. 4-264453 and No. 5-346682 disclose use of a hydrophobic-treated inorganic fine powder and a silicone-oil-treated inorganic fine powder in combination. Such methods are known in the art.\nMethods in which conductive fine particles are externally added as the external additive are also proposed in a large number. For example, carbon black as conductive fine particles is known to be used as an external additive for the purpose of providing conductivity to toners or controlling any excess charging of toners to make their triboelectric distribution uniform. Also, Japanese Patent Applications Laid-Open No. 57-151952, No. 59-168458 and No. 60-69660 disclose external addition of conductive fine particles such as tin oxide, zinc oxide and titanium oxide, respectively, to high-resistance magnetic toner particles. Japanese Patent Applications Laid-Open No. 61-275864, No. 62-258472, No. 61-141452 and No. 2-120865 disclose addition of graphite, magnetite, polypyrrole conductive particles or polyaniline conductive particles to toners.\nThese proposals, however, have room for further improvement to solve the above problems, when toner particles having small particle diameters are used in order to achieve a higher resolution.\nIn recent years, as copying machines and printers are being made compact, it has also become an important subject to achieve space saving, cost reduction and low power consumption. With regard to fixing assemblies, too, they have become required to be compact, structurally simple, and small power consumption. With this trend, toners are made to have a low viscosity at the time of melting to enlarge the area for their adhesion to fixing base materials or toner particles are incorporated with a release agent so that the toners can exhibit a sufficient fixing performance at a low amount of heat and a low pressure. Accordingly, binder resins used are required to have a low glass transition point (Tg) and a low molecular weight. However, for toners composed chiefly of soft components, it is difficult to achieve both fixing performance and high-temperature anti-offset properties simultaneously. Such toners also have a problem that they tend to cause a lowering of developing performance during long-term service or to stick or cling to the photosensitive member.\nMeanwhile, with regard to the improvement of fixing performance, various proposals have been made from old times. For example, Japanese Patent Publication No. 51-23354 discloses a pulverization toner improved in high-temperature anti-offset properties and low-temperature fixing performance, obtained by polymerizing a monomer such as styrene in the presence of a cross-linking agent and a molecular-weight modifier to obtain an appropriately cross-linked resin, and kneading this resin and a colorant such as carbon black, followed by pulverization. Japanese Patent No. 2681791 discloses a pulverization toner obtained by melt-kneading a styrene type binder resin containing a THF-insoluble matter in an mount of 10 to 60% by weight based on the weight of the resin, together with a charge control agent and a wax, followed by pulverization. These publications teach that molecular chains of the THF-insoluble matter (cross-linked component) of the binder resin are cut by melt kneading to form a high-molecular weight component, to thereby obtain a toner improved in both high-temperature anti-offset properties and low-temperature fixing performance. However, as a result of such thermal and mechanical cutting of the insoluble matter of the binder resin, soluble components formed by the cutting of molecular chains may have a fairly broad molecular weight distribution. Hence, medium-molecular-weight components which may damage low-temperature fixing performance also tend to be formed in a large quantity.\nSuch pulverization toner particles further have a problem that they have so low a circularity as to have a low transfer efficiency. Also, since the magnetic fine particles stand uncovered to toner particle surfaces in a large number, the toners tends to have a low fluidity and a low uniform chargeability.\nOn the other hand, in the case of the synthetic toners, different from the pulverization toners, toner particles can directly be produced without the step of melt kneading, and hence molecular chains of the insoluble matter (cross-linked component) formed at the time of polymerization are by no means cut. Thus, they are advantageous in that toner particles having very high anti-offset properties can be obtained, but on the other hand the insoluble matter tends to damage the low-temperature fixing performance. Accordingly, the low-temperature fixing performance and the high-temperature anti-offset properties must be balanced by controlling the insoluble matter. Also, when the magnetic fine particles are made insufficiently hydrophobic, the magnetic fine particles tend to stand uncovered to the toner particle surfaces in a large number, tending to make fixing performance poor and cause a deterioration of fixing assemblies.\nJapanese Patent Application Laid-Open No. 11-38678 discloses non-magnetic synthetic toner particles having 0 to 20% of components having a molecular weight of 1,000,000 or more and 0 to 60% of THF-insoluble matter, the total of the both being 1 to 60%. This publication, however, discloses a technique concerning non-magnetic toner particles, and there is room for improvement in respect of magnetic synthetic toner particles containing magnetic fine particles. In addition, Japanese Patent No. 2749234 discloses a process for producing magnetic toner particles in which a wax component in toner particles is present in a fibrous form. As disclosed in this publication, a polymerizable cross-linking agent is added in a monomer composition containing magnetic particles, which is then polymerized in the presence of an azo type polymerization initiator to obtain magnetic synthetic toner particles. Also, Japanese Patent No. 2749122 discloses a method of surface-treating magnetic particles with a polymer having a specific reactive group. As disclosed in this publication, a polymerizable cross-linking agent is added in a monomer composition containing magnetic particles, which is then polymerized in the presence of an azo type polymerization initiator to obtain magnetic synthetic toner particles. However, as presumed from the amount of the cross-linking agent, the type and amount of the polymerization initiator and the polymerization temperature which are described in these publications, the medium molecular weight component produced because of any excess formation of THF-insoluble matter or its cross-linking in a very weak state may be in a large proportion. Hence, in the case of the magnetic toner containing the magnetic fine particles in a large quantity, there is a problem on fixing performance. Also, in the magnetic synthetic toner particles obtained by the processes disclosed in these publications, the hydrophobic treatment of the magnetic fine particles used is insufficient, and there are problems also on fluidity and charging performance. Moreover, the achievement of both developing performance and fixing performance is also insufficient.\nWith regard to the image-forming method, as methods by which the electrostatic latent image is formed into a visible image, developing systems such as cascade development, magnetic brush development and pressure development are known in the art. Another method is also known in which, using a magnetic toner and using a rotary sleeve internally provided with a magnet, the magnetic toner is caused to fly across a photosensitive member and a developing sleeve by the aid of an electric field. For example, Japanese Patent Application Laid-Open No. 54-43027 discloses a method in which a magnetic toner is thinly coated on a magnetic toner-carrying member and this is triboelectrically charged, which is then served to develop an electrostatic latent image under application of a magnetic field. According to this method, the thin coating of a magnetic toner on a magnetic toner-carrying member enables the magnetic toner to be sufficiently triboelectrically charged. Moreover, the electrostatic latent image is developed while the magnetic toner is supported by the action of magnetic force. Hence, the magnetic toner can be kept from spreading to non-image areas, so that any fog can be kept from occurring and highly minute images can be obtained. Also, with regard to transfer efficiency, use of a toner having a uniform charge quantity distribution brings about a high transfer efficiency, but it is sought to make further improvement.\nSpherical toner particles are esteemed to have a high transfer efficiency. Concerning such particles, Japanese Patent Application Laid-Open No. 61-279864 discloses a proposal on a toner whose shape factors SF-1 and SF-2 are specified. Japanese Patent Application Laid-Open No. 63-235953 discloses a proposal on a magnetic toner made spherical by the action of mechanical impact force. However, toners are sought to be more improved in transfer efficiency.\nSuch spherical toner particles have on the one hand an advantage that they have a higher transfer efficiency than toner particles produced by pulverization, but on the other hand have a nature that they can be removed by cleaning with difficulty because of their sphericity. Moreover, since toner particles trend toward smaller particle diameters as state previously, toner particles may escape at the time of cleaning, and it has become more difficult to remove transfer residual toner completely by cleaning. However, an improvement of cleaning assemblies can keep the toner particles from escaping to a level that may cause no great problem. In image-forming methods having a conventional corona charging system, images having no problem in practical use can be formed.\nHowever, in recent years, from the viewpoint of environmental protection, in place of the primary charging and transfer process which have utilized corona discharging conventionally used, it is becoming prevailing to employ primary charging (contact charging) and transfer process (contact transfer) each of which make use of a member brought into contact with the photosensitive member surface, having great advantages of low ozone and low power consumption. For example, Japanese Patent Applications Laid-Open No. 63-149669 and No. 2-123385 disclose processes concerning the contact charging process and contact transfer process. In these processes, a conductive flexible charging roller is brought into contact with a photosensitive member and the photosensitive member is uniformly charged applying a voltage to the conductive roller, followed by exposure and development to form a toner image. Thereafter, another conductive roller to which a voltage is kept applied is pressed against the photosensitive member, during which a transfer medium is passed between them, and the toner image held on the photosensitive member is transferred to the transfer medium, followed by the step of fixing to obtain a fixed copy image.\nHowever, in such a contact charging process and a contact transfer process, too, there is room for further improvement. Stated specifically, in the case of the contact charging, the charging member is kept in pressure contact with the surface of the photosensitive member by pressing the former against the latter. Hence, the presence of any transfer residual toner tends to lower the contact between the contact charging member and the photosensitive member to tend to lower charging performance. In reverse development, the toner tends to spread to non-image areas to tend to cause fog. Also, any accumulation of toner on the charging performance tends to make it difficult to charge the photosensitive member uniformly, tending to cause a decrease in image density or cause coarse images. In addition, since the charging member is kept in pressure contact, melt-adhesion of toner tends to occur. These tendencies appear more remarkably as the transfer residual toner is in a large quantity.\nThen, in the case of the contact transfer, the transfer member is brought into contact with the photosensitive member through the transfer medium at the time of transfer, and hence the toner image is pressed when the toner image formed on the photosensitive member is transferred to the transfer medium, tending to cause a problem of partial faulty transfer, which is called xe2x80x9cblank areas caused by poor transferxe2x80x9d. Moreover, as a trend of techniques in recent years, there is a demand for developing systems of higher resolution and higher minuteness. To meet such a demand, toners are directed to have a smaller particle diameter. However, as toners are made to have a smaller particle diameter, the attraction force (e.g., mirror force or van der Waals force) of toner particles on the photosensitive member may increase to tend to result in an increase in the transfer residual toner, tending to cause faulty transfer.\nThus, in the image-forming methods making use of the contact charging process and contact transfer process which are very preferable taking account of environment, it is sought to bring forth a magnetic toner, and an image-forming method, which promise a high transfer performance and a superior charging stability and may hardly cause melt-adhesion of toner.\nMeanwhile, with regard to the toner having a high transfer efficiency as stated above, also proposed is a technique called a development-cleaning (also called cleaning-at-development) system or cleanerless system in which development and cleaning are carried out in the same step.\nDisclosure of conventional techniques concerning the development-cleaning or cleanerless system is, as seen in Japanese Patent Application Laid-Open No. 5-2287, focused on positive memory or negative memory appearing on images because of an influence of the transfer residual toner. However, in these days where electrophotography is utilized on and on, it has become necessary to transfer toner images to various recording mediums. In this sense, it is sought to make further adaptation to various recording mediums.\nThe prior art having disclosed the cleanerless system is seen in Japanese Patent Applications Laid-Open No. 59-133573, No. 62-203182, No. 63-133179, No. 64-20587, No. 2-302772, No. 5-2289, No. 5-53482 and No. 5-61383. These, however, neither mention any desirable image-forming methods nor refer to how the toner be constituted.\nAs developing systems in which the development-cleaning system or cleanerless system is preferably applied, in conventional development-cleaning systems basically having no cleaning assembly, it has been considered essential for the system to be so made up that the photosensitive member surface is rubbed with the toner and toner-carrying member. Accordingly, studies have largely made on contact developing systems in which the toner or toner-carrying member comes into contact with an image-bearing member. This is because, in order to collect the transfer residual toner in a developing means, it is considered advantageous for the system to be so made up that the toner or toner-carrying member comes into contact with and rub the image-bearing member. However, in the development-cleaning system or cleanerless system making use of a contact development system, its long-term service tends to cause deterioration of toner, deterioration of toner-carrying member surface and deterioration or wear of photosensitive member surface, but any satisfactory solution has not been made for running performance. Accordingly, it is sought to provide a development-cleaning system according to a non-contact developing system.\nHere, consider an instance in which the contact developing system is applied to a image-forming method employing the development-cleaning system or cleanerless system. In the image-forming method employing the development-cleaning system or cleanerless system, any cleaning member is provided and hence the transfer residual toner remaining on the photosensitive member surface comes into contact with the contact charging member as it is, to come to adhere to or mix in the contact charging member. Also, in the case of a charging system predominantly governed by a discharge charging mechanism, the transfer residual toner tends to adhere to the charging member because of a deterioration due to discharge energy. When insulating toners commonly used adhere to or mix in the contact charging member, the charging performance tends to lower.\nIn the case of the charging system predominantly governed by a discharge charging mechanism, the charging performance of the member to be charged tends to lower abruptly around the time when a toner layer having adhered to the contact charging member surface comes to have a resistance which may obstruct the discharge voltage. On the other hand, in the case of a charging system predominantly governed by a direct injection charging mechanism, the charging performance of the member to be charged may lower where the transfer residual toner having adhered or mixed has lowered the probability of contact between the contact charging member surface and the member to be charged.\nThis lowering of uniform charging performance of the member to be charged appears as a lowering of contrast and uniformity of electrostatic latent images after imagewise exposure to tend to cause a decrease in image density or make fog occur seriously.\nIn the image-forming method employing the development-cleaning system or cleanerless system, the point is that the charge polarity and charge quantity of the transfer residual toner on the photosensitive member is controlled so that the transfer residual toner can stably be collected in the step of development and the collected toner may not make the developing performance poor. Accordingly, the charge polarity and charge quantity of the transfer residual toner on the photosensitive member is controlled by means of the charging performance.\nThis will be described specifically taking the case of a commonly available laser beam printer. In the case of reverse development making use of a charging member for applying a voltage with negative polarity, a negatively chargeable photosensitive member and a negatively chargeable toner, in the transfer step the toner image is transferred to the recording medium by means of a positively chargeable transfer member. The charge polarity of the transfer residual toner varies from positive to negative because of its relation to the type of the recording medium (differences in thickness, resistance, dielectric constant and so forth) and the areas of images. However, the charging member having a negative polarity, used to charge the negatively chargeable photosensitive member, can uniformly adjust the charge polarity to the negative side even if even the polarity of the transfer residual toner has been shifted to the positive side in the transfer step. Hence, when the reversal development is employed as the developing system, the transfer residual toner, which stands negatively charged, remains at light-area potential areas to be developed by toner. At dark-area potential areass not to be developed by toner, the toner is attracted toward the toner carrying member in relation to the development electric field and is collected without remaining on the photosensitive member having a dark-area potential. That is, the development-cleaning system can be established by controlling the charge polarity of transfer residual toner simultaneously with the charging of the photosensitive member by means of the charging member.\nHowever, where the transfer residual toner has adhered to or mixed in the contact charging member beyond the contact charging member\"\"s capacity to control toner\"\"s charge polarity, it becomes difficult to uniformly adjust the charge polarity of the transfer residual toner. Also, even where the transfer residual toner has been collected on the toner-carrying member by mechanical force such as rubbing, the transfer residual toner may adversely affect the charging performance of toner on the toner-carrying member, resulting in a lowering of developing performance, unless its charge has not uniformly been adjusted.\nMore specifically, in the image-forming method employing the development-cleaning system or cleanerless system, the charge control performance at the time the transfer residual toner passes the charging member and the manner in which the transfer residual toner adheres to or mixes in the charging member are closely concerned with the running performance and image quality characteristics.\nIn order to prevent uneven charging to effect stable and uniform charging, the contact charging member may be coated with a powder on its surface coming into contact with the surface of the member to be charged. Such constitution is disclosed in Japanese Patent Publication No. 7-99442.\nThe contact charging member (charging roller) is follow-up rotated as the member to be charged (photosensitive member) is rotated (without no velocity differential drive), and hence may remarkably less cause ozone products compared with corona charging assemblies such as Scorotron. However, the principle of charging is still chiefly the discharge charging mechanism like the case of the roller charging mentioned previously. In particular, a voltage formed by superimposing AC voltage on DC voltage is applied in order to attain more stable charging uniformity, and hence the ozone products caused by discharging may more greatly occur. Accordingly, when the apparatus is used over a long period of time, difficulties such as smeared images due to ozone products tend to come out. Moreover, when applied in cleanerless image-forming apparatus, any inclusion of the transfer residual toner makes it difficult for the powder coated, to stand adhered uniformly to the charging member, so that the effect of carrying out uniform charging may lower.\nJapanese Patent Application Laid-Open No. 5-150539 also discloses that, in an image-forming method making use of contact charging, at least image-developing particles and conductive fine particles having an average particle diameter smaller than that of the image-developing particles are contained in a toner in order to prevent any charging obstruction which may be caused when toner particles or silica particles having not completely be removed by a cleaning means such as a cleaning blade come to adhere to and accumulate on the surface of the charging means during repetition of image formation for a long time. However, the contact charging used here, or proximity charging, applies the discharge charging mechanism, which is not the direct injection charging mechanism, and has the above problem ascribable to the discharge charging. Moreover, when applied in the cleanerless image-forming apparatus, nothing is taking into consideration about any of the influence on charging performance that is exercised when the conductive fine particles and transfer residual toner pass the charging step in a larger quantity than the apparatus having a cleaning mechanism, the influence on the collection of these large-quantity conductive fine particles and transfer residual toner in the developing step, and the influence on toner\"\"s developing performance that is exercised by the conductive fine particles and transfer residual toner thus collected. Furthermore, when the direct injection charging mechanism is applied in the contact charging, the conductive fine particles can not be fed to the contact charging member in necessary quantity to tend to cause faulty charging due to the influence of the transfer residual toner.\nIn the proximity charging, it is also difficult to uniformly charge the photosensitive member because of the large-quantity conductive fine particles and transfer residual toner, and the effect of leveling patterns of the transfer residual toner can not be obtained, to cause pattern ghost because the transfer residual toner may shut out pattern-imagewise exposure light. In-machine contamination due to toner may further occur when a power source is instantaneously put off or paper jam occurs during image formation.\nIn the image-forming method employing the development-cleaning system, development-cleaning performance can be improved by improving charge control performance required when the transfer residual toner passes the charging member. As a proposal therefor, Japanese Patent Application Laid-Open No. 11-15206 discloses an image-forming method making use of a toner having toner particles containing specific carbon black and a specific azo type iron compound and having inorganic fine powder. It is also proposed, in the image-forming method employing the development-cleaning system, to improve development-cleaning performance by reducing the quantity of transfer residual toner, using a toner having a superior transfer efficiency the shape factors of which have been specified. However, the contact charging used here also applies the discharge charging mechanism, which is not the direct injection charging mechanism, and has the above problem ascribable to the discharge charging. Moreover, these proposals may be effective for keeping the charging performance of the contact charging member from lowering because of the transfer residual toner, but can not be expected to be effective for positively improving the charging performance.\nIn addition, among commercially available electrophotographic printers, image-forming apparatus are also available which are designed for carrying out the development-cleaning system, in which a roller member coming into contact with the photosensitive member is provided between the transfer step and the charging step so that the performance of collecting the transfer residual toner at development can be assisted or controlled. Such image-forming apparatus have good development-cleaning performance and the waste toner can greatly be reduced, but involve a high cost and may damage the advantage inherent in the development-cleaning system also in view of compact construction.\nAs countermeasures for these, Japanese Patent Application Laid-Open No. 10-307456 discloses an image-forming apparatus in which a toner containing conductive charge-accelerating particles having particle diameter which is xc2xd or smaller than average particle diameter of toner particles or toner is applied in an image-forming method employing the development-cleaning system making use of the direct injection charging mechanism. According to this proposal, an image-forming apparatus for carrying out the development-cleaning system can be obtained, which can greatly reduce the quantity of waste toner and is advantageous for making the apparatus compact at a low cost, and good images are obtainable without causing any faulty charging and any shut-out or dispersion of imagewise exposure light.\nJapanese Patent Application Laid-Open No. 10-307421 also discloses an image-forming apparatus in which a toner containing conductive particles having particle diameter which is {fraction (1/50)} to xc2xd of average particle diameter of the toner is applied in an image-forming method employing the development-cleaning system making use of the direct injection charging mechanism and the conductive particles are made to have a transfer accelerating effect. Japanese Patent Application Laid-Open No. 10-307455 still also discloses that, a conductive fine powder is controlled to have particle diameter not larger than the size of one pixel of constituent pixels, and the conductive fine powder is controlled to have particle diameter of from 10 nm to 50 xcexcm in order to attain better charging uniformity.\nJapanese Patent Application Laid-Open No. 10-307457 discloses that, taking account of human visual sensation, conductive fine particles are controlled to have particle diameter of about 5 xcexcm or smaller, and preferably from 20 nm to 5 xcexcm, in order to make any influence of faulty transfer on images visually recognizable with difficulty.\nJapanese Patent Application Laid-Open No. 10-307458 also discloses an image-forming method which employs the development-cleaning system making use of the direct injection charging mechanism and in which a conductive fine powder is controlled to have particle diameter not larger than the average particle diameter of a toner to thereby prevent the conductive fine powder from obstructing the behavior of the toner at the time of development or prevent development bias from leaking through the conductive fine powder, and the conductive fine is controlled to have particle diameter larger than 0.1 xcexcm to thereby eliminate a difficulty that the conductive fine powder may become buried in the image-bearing member to shut out imagewise exposure light, thus superior image recording can be materialized.\nJapanese Patent Application Laid-Open No. 10-307456 discloses an image-forming apparatus which carries out the development-cleaning system and in which a conductive fine powder is externally added to toner particles so that the conductive fine powder contained in the toner particles may adhere to an image-bearing member in the step of development, at least at a contact zone between a flexible contact charging member and the image-bearing member, and may remain and be carried on the image-bearing member also after the step of transfer so as to stand between them, to thereby obtain good images without causing neither faulty charging nor shut-off of imagewise exposure light.\nIn all these proposals, however, there is room for further improvement in stable performances required when the apparatus are repeatedly used over a long period of time and in performances required when toner particles having a small particle diameter are used in order to achieve a higher resolution.\nAn object of the present invention is to provide a magnetic toner having solved the problems the prior art has had, and an image-forming method making use of the magnetic toner.\nAnother object of the present invention is to provide a magnetic toner having a good fixing performance, having superior environmental stability and charging stability and can form images in a high density and a high minuteness even over long-term service, and an image-forming method making use of such a magnetic toner.\nStill another object of the present invention is to provide an image-forming method which can well carry out the development-cleaning system.\nA further object of the present invention is to provide an image-forming method which can stably achieve a good charging performance and enables image formation by the cleanerless system.\nTo achieve the above objects the present invention provides a magnetic toner comprising magnetic toner particles containing at least a binder resin, a magnetic material containing a magnetic iron oxide, and a release agent.\nthe magnetic toner having;\na weight-average particle diameter of from 3 xcexcm to 10 xcexcm;\na magnetization intensity (saturation magnetization) of from 10 Am2/kg to 50 Am2/kg (emu/g) under application of a magnetic field of 79.6 kA/m (1,000 oersteds);\nan average circularity of 0.970 or more;\na ratio of weight-average particle diameter to number-average particle diameter, of 1.40 or less;\niron and an iron compound which stand liberated from the magnetic toner particles at a liberation percentage of from 0.05% to 3.00%; and\na resin component having a tetrahydrofuran (THF)-insoluble matter in an amount of from 3% by weight to 60% by weight.\nThe present invention also provides an image-forming method comprising;\na charging step of charging an image-bearing member electrostatically by applying a voltage to a charging member kept in contact with the image-bearing member, forming a contact zone between them;\nan electrostatic latent image forming step of forming an electrostatic latent image on the charged surface of the image-bearing member;\na developing step of forming a toner image by developing the electrostatic latent image by causing a magnetic toner to move to the electrostatic latent image at a developing zone where an alternating electric field is kept formed; the developing zone being formed between the image-bearing member for holding thereon the electrostatic latent image and a toner-carrying member for carrying the magnetic toner on its surface which are face to face disposed leaving a preset space between them, and a layer of the magnetic toner being formed on the surface of the toner-carrying member in a thickness smaller than that space; and\na transfer step of transferring the toner image to a transfer material via, or not via, an intermediate transfer member;\nthe steps being repeated to form images;\nwherein the magnetic toner comprises magnetic toner particles containing at least a binder resin, a magnetic material containing a magnetic iron oxide, and a release agent.\nthe magnetic toner having;\na weight-average particle diameter of from 3 xcexcm to 10 xcexcm;\na magnetization intensity (saturation magnetization) of from 10 Am2/kg to 50 Am2/kg (emu/g) under application of a magnetic field of 79.6 kA/m (1,000 oersteds);\nan average circularity of 0.970 or more;\na ratio of weight-average particle diameter to number-average particle diameter, of 1.40 or less;\niron and an iron compound which stand liberated from the magnetic toner particles at a liberation percentage of from 0.05% to 3.00%; and\na resin component having a tetrahydrofuran (THF)-insoluble matter in an amount of from 3% by weight to 60% by weight."}
{"text": "1. Technical Field\nThis is related to systems for separating flowing fluid into separate portions, and in particular, to systems for separating solid particles or gases from a liquid or liquids of differing buoyancies.\n2. Related Technology\nThe separation of mixtures into their component parts for the purposes of clarification, cleaning, and removal of undesirable portions of the mixture is a major function in many industries. In many cases it is desired to perform this function rapidly, at high flow rates and in an “on-line” operating mode. In many applications this function is performed by conventional centrifuges, often augmented with a settling or decanting process to overcome the inherent limitations of current centrifuges to achieve the desired level of end product purity needed.\nIn addition, the removal of air or other gases from fluid flows is widely applicable in cases in which the entrained gases may interfere with the function and performance of equipment. One prevalent application is in the removal of air from liquid coolant to improve the heat exchange performance of radiators. Another such application is the removal of air from lubricating fluid to improve the lubrication and cooling performance of the fluid. Yet another application is the removal of air from hydraulic fluid to improve the force transfer of hydraulic actuators. In all of these applications, the separated air or gas is typically removed, since recombination would tend to defeat the purpose of the removal. In these cases, air removal is the primary function and any suspended particles in the fluid are of little or no interest.\nIn other cases, recombination after separation may be desirable, for example, to retain oil fumes that are entrained in the airflow in order to preserve the total oil content of the system. In certain particle monitoring applications, such as that disclosed by U.S. Pat. No. 4,282,016 to Tauber et al., a goal is to remove particles above a certain size for analysis.\nReal time suspended particle monitors are disclosed in U.S. Pat. No. 5,572,320, U.S. Pat. No. 6,049,381, U.S. Pat. No. 7,921,739, and U.S. Pat. No. 8,056,400."}
{"text": "Acidic and alkaline cleaning compositions for hard surfaces have been used for many years to remove stubborn soils from a variety of surfaces found in household and institutional locations. Such soils include inorganic soils and soils derived from organic sources, such as fats, oils, proteins and carbohydrates. Such soils when heated can form hard tenacious deposits on a variety of surfaces including ceramic, stainless steel, concrete, tile and metal food preparation surfaces. Typical inorganic solids comprise insoluble materials derived from the hardness components of service water including substantially insoluble salts of calcium, magnesium, iron, manganese, etc. Such inorganic salts can be combined in some cases with organic residues which can form large, difficult to remove soil deposits. These soils can be unsightly, can take the form of large deposits of charred or baked-on residue, or large areas of white insoluble soap scum or hardness deposits. These soils can also promote or support the growth of microorganisms that can under certain circumstances contaminate food or other contaminatable materials or surfaces.\nA variety of cleaning compositions have been developed to deal with the tenacious organic and organic/inorganic matrix soils common in a variety of surfaces. One particularly useful form of cleaner is an aqueous alkaline cleaner commonly delivered from a pressurized aerosol or pump spray device. These types of cleaners have great utility for a variety of surfaces because the material can be delivered by spray to vertical, overhead or inclined surfaces or to surfaces having a complex curved or convoluted surface while achieving substantially complete coverage of the surface with the spray-on liquid cleaner. Acid spray-on cleaners are also known for removing basic inorganic soils and are becoming more common.\nOne substantial problem that arises with such spray-on materials relates to an airborne mist or finely divided aerosol generated during the spraying process. The aqueous liquid is converted by the action of propellant or pump action into a spray that is accompanied by a finely divided aerosol or mist. The spray portion contacts and remains on the target surface, while a substantial proportion of the aerosol can remain suspended in the atmosphere.\nSuch aqueous compositions having a strong acid or strong base cleaning component in the form of a finely divided aerosol or mist can cause respiratory distress in the user. Upon breathing the finely divided aerosol or mist, a very strong and irrepressible choking response is seen in most individuals that come in contact with irritating proportions of the aerosol produced by typical spray-on cleaners. The choking response is inconvenient, reduces cleaning efficiency in a variety of applications and in sensitive individuals can cause asthma attacks, respiratory damage, or other discomfort or injury.\nTo alleviate, to some degree, the choking response, some products have been formulated with reduced quantities of the strongly alkaline or strongly acid cleaning components to reduce the choking response. Strong caustic has been replaced by reduced alkalinity bases such as bicarbonate or by solvent materials. The reduction in concentration or substitution of these materials can often reduce the cleaning activity and effectiveness of the material when used.\nCrotty et al., U.S. Pat. No. 3,644,210, teach an alkaline cleaning material useful for removing burnt-on, baked-on food and grease from cooking surfaces comprising substantial proportions of alkali metal hydroxide in an aqueous cleaning base. Eisen, U.S. Pat. No. 3,779,933, teaches an aqueous alkali metal cleaning composition using clay and organic thickening materials. Mukai et al., U.S. Pat. No. 3,813,343, teach oven cleaning compositions containing an amine or ammonia base combined with dimethyl sulfoxide solvent and other compositions. Wise et al., U.S. Pat. No. 3,829,387, teach caustic cleaning compositions containing substantial proportions of sodium hydroxide, clay thickeners, and solvents in an aqueous base. Dillarstone et al., U.S. Pat. No. 4,214,915, and Canadian Patent No. 1,523,491, teach oven cleaning compositions using a relatively mild base such as sodium carbonate and sodium bicarbonate in an aqueous cleaning material. Culshaw, U.S. Pat. No. 4,676,920, teaches a thickened viscous scouring material using surfactants, abrasives and solvents for soil removal. De Buzzaccarini, U.S. Pat. No. 4,767,563, teach liquid scouring cleaning compositions using solvents, abrasives and surfactants for soil removal from hard surfaces. Cockrell, Jr., U.S. Pat. No. 4,877,691 (International Application No. PCT/US91/05092) teaches a barrier coating composition used to pretreat an oven to promote the subsequent removal of hardened baked-on soils that form on the pretreatment material. Dimond et al., Canadian Patent No. 1,047,903, teach oven cleaning compositions having reduced proportions of sodium hydroxide which are thickened using a bentonite in a substantially aqueous base. Silvester, Canadian Patent No. 1,211,674, teaches an improved oven cleaner using sodium bicarbonate as an alkali, glycerol as a solvent, a clay thickener in an aqueous formulation for oven cleaning. Flannery, United Kingdom Patent No. 2,019,876, teaches an aerosol comprising sodium bicarbonate as an oven cleaner.\nThe prior art taken as a whole shows that skilled artisans in experimenting with improving acid and alkaline hard surface cleaners, in particular alkaline oven cleaners, have attempted to reduce concentrations of sodium hydroxide to reduce the choking response. Ammonia or an organic base has been used to replace sodium hydroxide in reduced alkaline cleaners. Other relatively mild inorganic basic materials, such as sodium carbonate or sodium bicarbonate have also been used to replace sodium hydroxide. A variety of organic surfactants or glycol, alkyl ether or dimethyl sulfoxide solvent materials have been used to enhance the detergent properties of the reduced alkaline materials.\nThe prior art taken as a whole has attempted to replace sodium hydroxide to reduce choking response but does not recognize that the choking properties of these cleaners is related to the median particle size of the aerosol. For this reason, no attempt has been made in the prior art to formulate to adjust the particle size of the aerosol to reduce respiratory distress or choking reflex."}
{"text": "In recent years, wireless power transfer techniques have been gaining attention in order to provide power supply or perform charging. Research and development are being conducted regarding a wireless power transfer system wirelessly performing power transfer to various electronic apparatuses such as mobile terminals and notebook computers and household electrical appliances or to power infrastructure equipment.\nIn order to use wireless power transfer, it is preferable to standardize so that no problem occurs in the use of a power source of a power source and a power receiver of a power receiver that are of different manufactures.\nAmong conventional wireless power transfer techniques, a technique using electromagnetic induction and a technique using radio waves have generally been known. On the other hand, expectations for power transfer techniques using magnetic field resonance (magnetic resonance) or electric field resonance have been increasing recently, as techniques allowing for power transfer to a plurality of power receivers and power transfer to various three-dimensional postures while maintaining some distance between power sources and the power receivers. Electric field resonance may also be called electric resonance.\nAs described above, attention has conventionally been paid to wireless power transfer techniques for wirelessly transferring power for the purposes of power supply or charging. Nevertheless, standardization of power transfer techniques, for example, using magnetic field resonance or electric field resonance has not been made so far.\nThere has thus been a concern over stagnation of practical application of a power transfer system using magnetic field resonance or electric filed resonance or of a power source and a power receiver.\nA variety of wireless power transfer techniques have conventionally been proposed.    Patent document 1: Japanese Laid-open Patent Publication No. 2010-239769    Patent document 2: U.S. Pat. No. 7,825,543    Non-Patent document 1: SHOKI Hiroki, et al., “Standardization Trends on Wireless Power transfer”, Technical Report of The Institute of Electronics Information, and Communication Engineers (IEICE technical report), WPT 2011-19, December 2011."}
{"text": "A video capture apparatus can record a moving image as a series of still frames. The series of still frames may be digitally recorded as image data. When outputted, the series of still frames are displayed at a consistent rate. The consistent rate allows a viewer to see the image captured by the apparatus as a moving picture. Before viewing the series of image frames, the image may be stored at a local media player. When stored, the series of frames can consume a relatively large amount of computer memory. Therefore, video compression can be used to conserve memory.\nIn video compression, algorithms are applied to the image data. The algorithms attempt to reduce an amount of storage space that the image data consumes. When the image data is displayed, a reverse algorithm is applied to the image data. The reverse algorithm allows a near identical image or an identical image of the retained image to be displayed. Many algorithms used in compression affect the quality of the displayed image and may be inefficient. A more advantageous frame and storage management system may be desired."}
{"text": "1. Field of the Invention\nThis invention relates to an in-vehicle electronic control device which is a constant voltage control device built into an in-vehicle electronic control device such as an engine control device and a transmission control device, the in-vehicle electronic control device being so improved as to reduce a fluctuation error in output voltage of a constant voltage control circuit generating a predetermined constant voltage output by the use of power supplied from an in-vehicle battery.\n2. Description of the Related Art\nIn the in-vehicle electronic control device, it is important to enhance a constant voltage control accuracy of the built-in constant voltage power source circuit in view of, for example, improvement in AD conversion accuracy in a multi-channel AD converter and improvements in generation accuracy of various comparison reference voltages.\nHowever, though an individual fluctuation of each component parts is inevitable, an expensive constant voltage power source circuit will undesirably be required for obtaining a very high accuracy constant voltage output.\nIn order to avoid such problem, the constant voltage control accuracy has been improved by using a constant voltage power source device having a constant voltage control accuracy lower than an expectation value and adding an auxiliary unit for correcting a reduction in constant voltage control accuracy.\nFor instance, Patent Document 1: Japanese Patent Publication 2002-366238 “Circuit Device and method for Setting Adjustment Data of Circuit Device” (refer to paragraph [0009] and FIG. 5) describes that in a constant voltage generation circuit included in a semiconductor sensor or an engine control unit (ECU) for a vehicle, a reference voltage VBGR which is an output voltage of a reference voltage generation circuit having a band gap regulator is amplified by an operational amplifier to obtain a predetermined constant voltage output Vcc, and, at the same time, a predetermined constant voltage output Vcc is obtained by writing adjustment data to an EPROM memory in order to correct an individual fluctuation in reference voltage VBGR and by changing a gain of the operational amplifier by the adjustment data.\nIn the invention disclosed in Patent Document 1, functions not only improving the constant voltage control accuracy but also adjusting a threshold value of a reset circuit in tandem with the adjustment data are added to the unit.\nAlso, Patent Document 2: Japanese Patent Publication 09-297623 “Voltage/Current Regulator Circuit” (refer to FIG. 3 and Abstract) describes that a resistance circuit network deciding a resistance value by the use of ROM data is included as a resistance circuit for deciding a differential amplification ratio in a constant voltage regulator circuit which is used for a temperature compensation oscillator, for example, and outputs signals subjected to differential amplification between a reference voltage source and a detected voltage, and fine regulation of a stabilized voltage is achieved by the use of the ROM data."}
{"text": "Extracorporeal shockwave therapy (herein referred to as ‘ESWT’) is non-surgical, non-invasive treatment of medical conditions using acoustic shockwaves. First use of shockwave therapy in the early 1980's was utilized to fragment kidney stones termed shockwave lithotripsy. Continued development of shockwave treatment showed the possibility of stimulating bone formation, angiogenesis, chronic orthopedic inflammation healing, bone healing (osteogenesis), wound healing, revascularization, angiogenesis are well known and described in medical literature.\nA shockwave is a form of acoustic energy resulting from phenomena that create a sudden intense change in pressure for example an explosion or lightning. The intense changes in pressure produce strong waves of energy that can travel through any elastic medium such as air, water, human soft tissue, or certain solid substances such as bone.\nAcoustic shockwaves may be generated by various methods, electrohydraulic (also referred to as spark gap), electromagnetic (also referred to as ‘EMSE’), piezoelectric and ballistic shockwave.\nEach method needs an apparatus to focus the generated shockwave so as to provide a focal point and/or focal zone for the treatment area. In the focal zone shockwaves produce much higher pressure impulses as compared with the zones outside of the focal zone.\nMechanical means for focusing each of these methods is generally realized with an appropriate arrangement of surfaces reflecting the wave toward the desired focal point and/or an appropriate arrangement of the generating devices.\nSpark gap systems incorporate an electrode (spark plug), to initiate a shockwave, and ellipsoid to focus the shockwave. EMSE systems utilize an electromagnetic coil and an opposing metal membrane. Piezoelectric systems form acoustical waves by mounting piezoelectric crystals to a spherical surface to provide focus. Of the three systems, the spark gap system is generally preferred in the art for generating therapeutic shockwaves ESWT as it introduces more of the generated shockwave energy to the treatment target site.\nIn spark gap systems, high energy shockwaves are generated when electricity is applied to an electrode positioned in an ellipsoid immersed in treated water. When the electrical charge is fired, a small amount of water is vaporized at the tip of the electrode and a shockwave is produced. The shockwave ricochets from the side of an ellipsoid and converges at a focal point, which may then be transferred to the area to be treated.\nIn electromagnetic systems an electrical impulse is circulated in a coil. The coil produces an electromagnetic field that expels a metallic membrane to produce the mechanical impulse.\nIn piezoelectric systems ceramic material with piezoelectric characteristics is subjected to an electrical impulse. The electric impulse modifies the dimension of the ceramic material to generate the desired mechanical impulse. A focal point is attained by covering a concave spherical surface with piezoelectric ceramics converging at the center of the sphere.\nThe electrohydraulic, electromagnetic, and piezoelectric are all forms of shockwave generators that utilize high voltage power sources from 10 kV to about 25 kV in order to generate the required shockwave of about 100 bar to about 1000 bar). The drawbacks for such high voltage shockwave technology includes limitations both relating to the actual treatment and to the actual device and system. Treatment related limitations for example include production of a limited focal zone treatment area and low treatment efficacy. System limitations for example include cost, size, and durability where systems are generally expensive, large, heavy and require frequent maintenance. However, the biggest limitation of such system relates to the operating costs where such systems require many disposable accessories and integral electronic parts.\nWhile high voltage device produce a shockwave pressure wave of about 100 bar to about 1000 bar, state of the art ballistic shockwave system offer generation of low level shockwaves, having pressure wave from about 50 bar to about 150 bar. As its name suggests ballistic shockwave system generate shockwaves as a result of a ballistic collision between a projectile and a generating surface. The projectile is accelerated and allowed to collide with the shockwave generating surface.\nState of the art, ballistic shockwave systems are utilized for medical applications such as in physiotherapy applications for example for the treatment of inflammations and/or in dermatology and cosmetic applications, for example in the treatment of cellulite.\nCurrent ballistic shockwave systems are limited in that a low pressure gas source (1-6 bar) leads to shockwaves that have low tissue penetration, small treatment and/or focal zone, high rates of re-treatment, discomfort due to the applicator's movement during the ballistic collision, are not readily mobile as they require an air compressor to produce the appropriate pressure. Other prior art ballistic shockwave system utilize an operational pressure of 15-30 bar, for example as described in U.S. Pat. No. 7,470,274 to Lebet. Moreover ballistic shockwave system generally do not provide for non-invasive extracorporeal shockwave treatment.\nSimilarly other prior art US2005/0209586, U.S. Pat. No. 6,413,230, WO2003084608, WO2008/007502 A1, WO2008/145273, U.S. Pat. No. 6,736,784, WO2010049519 describe ballistic shockwave systems using low level shock wave production by low gas pressure ballistic technology.\nOther forms of ballistic shockwave generators include electromagnetic ballistic systems are further limited in that they tend to heat up and therefore require a cooling system due to the inclusion of an electromagnetic components."}
{"text": "1. Field of the Invention\nThe current invention generally relates to digital logic systems. More specifically, the current invention relates to capacity on demand using control over latency and bandwidth of signal bussing.\n2. Description of the Related Art\nComputing systems are currently available that provide capacity on demand. Capacity is used to denote a performance characteristic of a computing system. For example, in a commercial workload, ability to handle some number of “transactions per minute” is one measure of capacity. Ability to provide a fast response time to a request is another measurement of capacity. Ability to perform some number of floating point operations per second is a third example of capacity. Demand is the workload being placed on the computing system. A computing system having a large capacity but only seeing a small demand is wasteful and expensive. A computing system having a demand higher than a capacity provides slow response to the user. Demand tends to vary greatly during a day or day of week (e.g., in many cases, during weekends, demand is low on many computing systems).\nIBM Corporation of Armonk, N.Y., currently provides capacity on demand (COD), which is sometimes called Capacity Upgrade on Demand (CUoD), for IBM eServer pSeries p650, 670 and 690 computing systems. Reference “pSeries Capacity Upgrade on Demand advantages” viewable at the following URL:\nhttp://www-8.ibm.com/servers/eserver/au/pseries/cuod/advantages.html\nThis capability allows a customer to accommodate unexpected demands on a computing system installed. For example, pSeries 670 and 690 servers are available in units of four active and four inactive processors with up to 50% of the system in standby. As workload demands require more processing power, unused processors can be activated simply by placing an order to activate the additional processors, sending current system configuration to an authorizing source and receiving over the internet an electronically encrypted activation key which unlocks the desired amount of processors. There is no hardware to ship and install, and no additional contract is required. Memory activation works the same way. CUOD is available in various sizes for the p650, p670 and p690 systems. Activation in 4 GB (Gigabyte) increments is made by ordering an activation key to unlock the desired amount of memory.\nProviding CUoD by enabling entire processors limits granularity of capacity upgrades (or reductions). For example, in a computing system having four active processors, the smallest increment in capacity is one processor.\nProviding CUoD by adding a processor, in some computing environments, may not provide a proportional increase in capacity. For example, if a computing system is running a numerically intensive program and an improvement in that numerically intensive program is desired, adding a processor will not provide the desired improvement, unless the numerically intensive program is capable of distribution across more than one processor. Adding a large increment of memory may not be of much help in the numerically intensive program, either, since many such programs make extensive use of a relatively small amount of memory, and bandwidth, rather than total memory size, is the dominant consideration.\nAdditional references include patent applications filed by the current assignee of the present patent application include: Ser. No. 10/616,676, “Apparatus and Method for Providing Metered Capacity of Computer Resources”, by Daniel C. Birkestrand et al, filed Jul. 10, 2003; Ser. No. 10/406,164, “Billing Information Authentication for On-Demand Resources”, by Daniel C. Birkestrand et al, filed Apr. 03, 2003; and Ser. No. 10/640,541, “Capacity On Demand Grace Period for Incompliant System Configurations”, by Daniel C. Birkestrand et al, filed Aug. 28, 2003.\nTherefore, there is a need for a method and apparatus to provide a finer granularity of capacity on demand."}
{"text": "With the continuing growth in performance of computing technology, the interconnect bandwidth required among microprocessors, memories, and input/output devices continues to increase. Such high speeds create increasing problems for the technology used for interconnections among the microprocessors, memories, and input/output devices. For example, copper trace technology, such as conventionally employed on printed circuit boards, is expected to be limited to 15-20 Gigabits per second as a result of signal degradation, power dissipation, and electromagnetic interference unavoidable at such high clock speeds.\nManufacturers of many electronic products have sought to address the limitations of copper by using more exotic substrates with lower dielectric loss or using more sophisticated input/output equalizers at the transmitter and receiver. Unfortunately all of these potential solutions are costly and power consumptive. Therefore traditional interconnect scaling will no longer satisfy performance requirements. Defining and finding solutions beyond copper and low dielectric loss material will require innovation in design, packaging and unconventional interconnect technology.\nOne alternative attracting increasing attention is the use of optical interconnect technology. Optical communication technology has been used for many years in long distance applications such as telephony and the internet, and is now sometimes implemented for use in shorter distance applications for the enterprise such as storage area networks and rack-to-rack interconnections. Such optical technology has already demonstrated that at high frequencies, the optical fibers provide longer distance-higher bandwidth capability as compared with electrical cables, yet minimize loss in the transmitted signal.\nUnfortunately, optical interconnect technology has associated with it a separate set of implementation difficulties. These difficulties include optical coupling efficiencies, fiber alignment, complex packaging technologies, including hermeticity, thermal management, electrical performance, and manufacturability—the ease of assembly and amenability to automation. Thus, generally speaking, current optical modules or transceivers like SFP, 300 pins MSA, Xenpack, Xpak, X2 and XFP form factors provided by various suppliers including Emcore, Finisar, Agilent, and Bookham are still complex to manufacture and expensive.\nAdditionally for substantial utility in the computing market, optical transceivers also need (1) extremely low power consumption to compete with electrical solutions, (2) a smaller form factor adapted for and appropriate to the computer industry, and (3) very large interconnect bandwidth to provide data consummate with processing power. One example of a prior art solution is the IBM Terabus project (Exploitation of optical interconnects in future server architectures, Benner et al., 49 IBM J. Res. & Dev. No. 4/5, July/September 2005). In this concept the transmitter and receiver modules are separate and require a complex structure with multiple chips/modules mounted on a substrate—one module with an array of transmitters and another module with an array of receivers.\nA novel solution adapted for computing application to interconnect processing units while increasing communication bandwidth, maintaining an ultra small form factor, decreasing power consumption, simplifying and enabling automated assembly is required."}
{"text": "Patient beds used in health care facilities and home care settings often include a lift system allowing a patient or caregiver to adjust the height of the bed. The lift system must satisfy a number of potentially conflicting constraints. For example, the lift system should be quiet, dependable, safe and damage resistant. It should also be inexpensive to manufacture and should be adaptable to different bed models with no more than simple, inexpensive modifications. Because the lift system typically resides underneath the elevateable components of the bed, it must be compact enough to allow the bed to be positioned at very low elevations and yet must also have enough reach to position the bed at elevations high enough to be satisfactory for the caregiver. Compactness also makes space available for other under-bed components."}
{"text": "1. Field of the Invention\nThe present invention relates to an adaptive writing method for a high-density optical recording apparatus and a circuit thereof, and more particularly, to an adaptive writing method for optimizing light power of a light source, e.g., a laser diode, to be suitable to characteristics of a recording apparatus, and a circuit thereof.\n2. Description of the Related Art\nWith the multi-media era requiring high-capacity recording media, optical recording systems employing high-capacity recording media, such as a magnetic optical disc drive (MODD) or a digital versatile disc random access memory (DVD-RAM) drive, have been widely used.\nAs the recoding density increases, such optical recording systems require optimal and high-precision states. In general, with an increase in recording density, temporal fluctuation (to be referred to as jitter, hereinafter) in a data domain increases. Thus, in order to attain high-density recording, it is very important to minimize the jitter.\nConventionally, a write pulse is formed as specified in the DVD-RAM format book shown in FIG. 1B, with respect to input NRZI (Non-Return to Zero Inversion) data having marks of 3T, 5T and 11T (T being the channel clock duration), as shown in FIG. 1A. Here, the NRZI data is divided into mark and space. The spaces are in an erase power level for overwriting. The waveform of a write pulse for marks equal to or longer than 3T mark, that is, 3T, 4T, . . . 11T and 14T is comprised of a first pulse, a last pulse and a multi-pulse train. Here, only the number of pulses in the multi-pulse train is varied depending on the magnitude of a mark.\nIn other words, the waveform of the write pulse is comprised of a combination of read power (FIG. 1C), peak power or write power (FIG. 1D) and bias power or erase power (FIG. 1E). Here, the respective power signals shown in FIGS. 1C, 1D and 1E are all low-active signals.\nThe waveform of the write pulse is the same as that in accordance with the first generation 2.6 GB DVD-RAM standard. In other words, in accordance with the 2.6 GB DVD-RAM standard, the waveform of the write pulse is comprised of a first pulse, a multi-pulse train and a last pulse. Although the rising edge of the first pulse or the falling edge of the last pulse can be read from a lead-in area to be used, adaptive writing is not possible since the write pulse is fixed to be constant.\nTherefore, when a write operation is performed by forming such a write pulse as shown in FIG. 1B, severe thermal interference may occur back and forth with respect to a mark in accordance with input NRZI data. In other words, when a mark is long and a space is short or vice versa, jitter is most severe. This is a major cause of lowered system performance. Also, this does not make it possible for the system to be applied to high-density DVD-RAMs, e.g., second generation 4.7 GB DVD-RAMs."}
{"text": "A video mosaic effect is created using a conventional digital video effects system by compressing the size of the image represented by an input video signal, thereby reducing its resolution, then progressively expanding the reduced-resolution image until a desired size is attained. At that point, the expanded, reduced-resolution image is replicated over the output scene, whereby the output scene is composed of multiple rectangular mosaic tiles. The nature of the mosaic effect is limited because the mosaic tiles are all identical."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a nail gun, and more particularly to a safety structure of a nail gun.\n2. Description of the Related Art\nTypically, to prevent a nail gun from being triggered unexpectedly to hurt somebody, the nail guns are equipped with a safety for protection of user and people around."}
{"text": "Prior to the initial experimentation with double-finned surfboards in the early 1970's, a single center fin, located at the very tail of the board, provided the directional stability essential to the basic performance of the board. Since the advent of tri-fin or “thruster” type surfboards in the early 1980's, high-performance surfboards have also incorporated two side-fins to dramatically increase the board's speed and maneuverability. The side-fins are located on opposite sides of the board near the perimeter edge or “rail,” and well forward of the single, central trailing fin at the tail.\nIn the tri-fin configuration, it is well established that the center fin is primarily a stabilizing fin and functions in a manner very similar to the fixed keel on a sailboat or the vertical stabilizer on an aircraft—i.e. if the board yaws or departs from its original heading, the rotation of the board causes the water-flow to strike the fin at an angle; this creates a low-pressure area on the opposite or lee side of the fin that resists the yaw, and allows directional stability to be maintained.\nKnowledge is still very limited, however, as to how the side-fins enhance the speed and maneuverability of modern multi-finned type boards. This has long been a major problem in surfboard design. As a result, the first, largely experimental “twin-fin” and “fish” style surfboards, the double-finned predecessors of the modern tri-fin, suffered for many years from a variety of poorly understood control problems. The early control problems—which were collectively referred to as “tracking”—were found to be greatly reduced by using a negatively angled side-fin setting. Although this eliminated the original tracking problem, it also caused an overly loose, drifting type of turn that many riders, even at the expert level, found very difficult to control. Eventually, the problem was remedied by adding a third stabilizing fin at the very tail of the board, the configuration in current use today. Though much faster and more maneuverable than the single-finned board types that preceded it, the current tri-fin setting was arrived at almost entirely through trial and error; as a consequence, it retains features that actually contribute to a marked increase in drag. The main drawbacks of prior art tri-fins may be summarized briefly as follows:\nEach side-fin is set at a negative angle of attack or “toe-in” angle of between three and five degrees, so that the leading edge points in the approximate direction of the longitudinal centerline at the nose. The angle is measured using the chord line (a straight line drawn through the leading and trailing edges of the fin at the fin base), which is referenced to the longitudinal centerline provided by the wooden center spar or “stringer” that runs the length of the board. The negative angle of attack or toe-in causes the water-flow to strike the side-fins at an angle, and creates high drag from the “snowplow” effect when the rider's weight is neutrally centered on the board.\nThe cambered foil of the side-fin adds to this drag: in the longitudinal cross-section view commonly used to depict the airfoil section of a wing, the foil of the side-fin is asymmetrical, and has an average curvature greater on one side than the other. The foil of the conventional prior art side-fin is flat to slightly concave on the inside surface (the side facing the longitudinal centerline or stringer), and curved on the outside (the side facing the perimeter edge or “rail”). Although the cambered side-fin foil appears to give better performance and greater average speed, knowledge is currently very limited as to the reasons why, since both the flat-sided, and particularly the slightly concave side-fin foil, would appear to greatly increase the drag from the negative toe-in angle. It is well known that separation of the boundary layer and turbulence occurs more readily when a flat or concave surface is set at an angle to a fluid flow, versus a symmetrical foil, for example, where both sides are convex and curve equally in opposite directions in a low-drag, streamlined shape.\nCurrently, the rider can overcome the high drag of the side-fin setting by constantly turning the board. As noted above, the high drag condition occurs primarily when the rider's weight is neutrally centered on the board—the drag is reduced, however, when the rider leans to initiate a turn and lifts the opposing side-fin free; the angle of the side-fin remaining in the water then acts like a deflected rudder and aids the board's rotation in the turn; on a tri-fin board, the rider's normal weight shift further in the turn will then set the center stabilizing fin, and prevent the overly loose, difficult to control, drifting type of turn that, subsequent to the “tracking” problem, was the major drawback that greatly limited the acceptance of the early double-finned style boards. Surfboard designers have long noted that adding a third stabilizing fin does little to diminish the maneuverability of the board—it instead produces such a noticeable burst of speed and acceleration in a turn that, in the early development of the tri-fin, the center stabilizing fin almost immediately came to be referred to as “thruster” fin, and the tri-fin set-up as a “thruster” type board. In the tri-fin or thruster configuration, however, the addition of the center stabilizing fin causes a third and final drawback:\n3) The location of the center stabilizing fin is precisely the opposite of the optimum theoretical configuration: i.e., if the negatively angled side-fin functions as a deflected rudder, it should be placed as far behind the board's axis of rotation as possible so as to increase its moment arm; the added leverage would lessen the surface area of the side-fin and the amount of negative toe-in angle required for a given turning moment, and thereby reduce drag. Locating the fin or fins required for directional stability forward of a negatively angled trailing fin, closer to the axis of rotation, would increase the directional instability of the fin-setting by allowing the negatively angled rearward fin to truly function as a permanently deflected rudder. Failure to correct the drawbacks outlined above, and the absence of innovation regarding fin placement on multi-fin type boards (the group includes other multi-finned variants, e.g., “twinzers,” “quads,” “fishes,” etc. all of which use the negatively angled side-fin setting), is largely due to the poor understanding of the role the fins play in enhancing the performance of the board. Despite the high speed and exceptional maneuverability of modern multi-finned boards vs. the early single-finned board types, at present, their higher performance actually comes at a cost of considerable drag. From a hydrodynamic standpoint, it can be seen that the board-making arts currently have need of a cambered side-fin foil that exhibits reduced drag at the conventional negatively angled side-fin setting, as well as multi-fin arrangements that will introduce directional instability, but at a reduction in drag over the multi-fin configurations of the prior art.\nThe following description is intended to impart an understanding of the present invention to a person skilled in the art of surfboard design. Those skilled in the art, however, will be aware of the current lack of tank-testing facilities, and the absence of any method that can accurately duplicate a breaking wave, the movement of the board on a wave, or the effects of the rider maneuvering the board in a controlled setting. Therefore, at least some of the material disclosed herein is a subjective interpretation of observed phenomena, and the descriptions provided below should not be interpreted in a way that will limit the invention, which is defined more fully and accurately in the appended claims.\nAt the time the present invention was made, the board-making arts lacked an explanation for the clearly superior performance of multi-finned type boards. As will be appreciated by those skilled in the art upon reviewing the disclosure below, the much higher speed of currently available multi-finned boards can be largely attributed to the higher lift coefficient of the cambered side-fin foil. The following detailed description of the invention therefore begins with a discussion of the relationship between the (hydro-) foil of the fin, and the airfoils of a wing and a sail, which respond in similar ways to a fluid flow despite the differing densities between air and water.\nSailboats and aircraft are able to maneuver because of the differential “lift” of a plurality of separate air- and hydrofoils at differing angles of attack: on a sailboat, for example, the “lift” of the deflected rudder creates a yawing moment behind the fixed keel that causes the sailboat to rotate in a turn; on an airplane, the differential lift between the wing and the horizontal tail (as altered by deflected control surfaces such as ailerons, elevons, the elevator, etc.) makes it possible for the aircraft to execute banked turns and fly in a loop. The board-designer, therefore, may use the same principles and analyze the angle of attack of the fin(s) relative to the direction of the water-flow through a turn, and arrange the fins, and the foil of the fins, to optimize the speed and performance of the multi-finned board as it is maneuvered on a wave.\nBoard designers may therefore benefit from a fuller knowledge of the similarities between the hydrofoil of the fin and the airfoil of the wing and sail, and make use of the extensive aeronautical research that has been compiled comparing the performance of various airfoil sections at different wind speeds and angles of attack. As shown in greater detail below, aeronautical engineers have developed sophisticated means of accurately measuring the performance of a wing; typically, the relevant wind tunnel data are plotted in graph form or, as shown in FIG. 1A and FIG. 1B, by using vectors, in which the length and direction of an arrow indicates the magnitude and direction of the force of the air pressure, or pressure field, that develops around the airfoil of a wing in response to its incidence, or angle of attack, relative to the airstream. For illustration purposes, the vectors shown in FIG. 1A and FIG. 1B, which actually represent the pressure differential around the airfoil of a wing, will be assumed to be completely interchangeable with the flat-sided cambered side-fin foil of the prior art. In addition, although the foils in FIG. 1A and FIG. 1B are depicted in a vertical orientation, in the following discussion they will be referred to as being in a horizontal position when the description is of an airfoil in flight, while the fluid flow F will be understood to represent both air- and water-flow.\nIn FIG. 1A, the vectors shown represent the pressure differential typically seen around the airfoil of a wing at cruise, when the airstream or airflow F is almost parallel to the airfoil of the wing. Ordinarily, the aircraft is designed so that the airplane's fuselage is completely level under normal flight conditions for minimum drag, while the wing is positioned at a very low but slightly positive angle of attack (e.g., typically about two degrees), so that the highest pressure will be at the leading edge of the wing, as shown, while the much lower pressure on the upper surface of the airfoil holds the aircraft aloft.\nIn aircraft design, a basic problem is that the pressure field depicted in FIG. 1A is unequal; as a result, the wing has a “pitching moment” and the aircraft tends to nose downward until the pressure around the wing is equalized. To prevent this, a horizontal stabilizer is provided at the tail, the airfoil of which is set at a slightly negative incidence or angle of attack so as to provide steady downward pressure, which counters the pitching moment of the wing and allows the aircraft to remain in steady, level flight.\nComparing the foil of a board fin to the airfoil of the wing, it can be assumed that a parallel side-fin setting will create a “yawing moment” similar to the pitching moment of the wing, and create control problems that would require a negatively angled trailing fin to counter, assuming the example set in aircraft design is followed. In surfboard design, however, the “tracking” problems exhibited by the very early fish style boards, which originally used a parallel side-fin setting, were eliminated by changing the fin position so the side-fin was fixed at a negative angle of attack. Despite the high drag and snowplow effect of the now standard, negatively angled side-fin setting, the modern multi-finned board type is much faster than the single-finned board types that preceded it. As will be appreciated by those of skill in the art after reading the disclosure below, this is because the rotation of the board in a turn places the side-fin foil at a high angle of attack, and a pressure differential forms around the fin that is much like the airfoil of a wing or sail at a similar angle of attack, as described in greater detail below.\nIn FIG. 1B, the pressure differential shown is typical of an airfoil at a very high angle of attack, when the airflow F is striking the underside of the wing, as is the case when the aircraft is flying in a loop or pulling out of a dive. Note that in either case the motion of the aircraft describes an arc, and that the direction of the airflow F is almost entirely due to the motion of the aircraft itself (assuming a still day with little breeze). When the airfoil is at a high angle of attack as shown, a very large area of negative pressure develops around the leading edge of the airfoil and pulls the wing forward. It is known that a similar area of low-pressure around the forward portion of a sail drives a sailboat forward and enables it to sail into the wind. From FIG. 1B, it can be assumed that if the rotation of the board through a turn places the fin at a correspondingly high angle of attack, an area of very low pressure will develop around the leading edge of the fin and accelerate the board forward; the aforementioned effect provides an explanation for the greater speed of multi-finned type boards.\nIn terms of board design, however, it is equally important to note that the pressure differential between the leading and trailing sections of the airfoil in FIG. 1B is very large; hence, an airfoil at a high angle of attack tends to have a very large pitching moment (in the case of a wing) or yawing moment (in the case of a sail), the effects of which must be countered with considerable deflection of the elevator or rudder to maintain directional control. It can be assumed that the cambered side-fin foil at a similarly high angle of attack will also have a very large yawing moment, and that the yawing moment will be opposite the rotation of the turn. The reverse yawing moment of the side-fin in a turn provides an effective explanation for the poorly understood control problems exhibited by the original wide-tailed twin-fins and the very early double-finned fish style boards of the prior art.\nAs previously discussed, the “tracking” problems of the original double-finned boards were eliminated through trial and error, without benefit of the information provided in the discussion above. As a consequence, current multi-fin configurations retain a number of features that actually contribute to a marked increase in drag. The source of the drag is illustrated in more detail in FIG. 2, which depicts the bottom of a conventional tri-fin surfboard according to the prior art. As shown, the two side-fins are located on opposite sides of the board near the perimeter edge or “rail,” and well forward the center stabilizing fin at the tail. When the board is at speed on the wave and the rider's weight is neutrally centered on the board, the heading H of the board will cause a water-flow F that is substantially opposite the heading; when the water-flow F parallels the longitudinal centerline or stringer as shown, the negatively angled side-fin setting, which has a standard toe-in angle of approximately four degrees, causes the water-flow F to strike the outside, cambered surface of the side-fins (the side facing the perimeter edge or rail), and creates high drag due to the low-pressure area (depicted here as turbulence) that develops on the lee or inside surface of the side-fins (the side facing the longitudinal centerline or stringer).\nFIG. 2A and FIG. 2B are closer, cross-section views depicting the cambered foil of prior art side-fins. The conventional flat-sided cambered foil of the prior art is shown in FIG. 2A; for a given thickness, the prior art foil shown in FIG. 2B has slightly increased camber due to the shallow concave of the inside surface. The views depict how the negative toe-in of the side-fin causes the water-flow F to strike the side-fins at an angle, which causes the water-flow on the lee or inside surface of the side-fins to tend to separate or become turbulent, and increases drag.\nNote that the actual angle of the side-fin foil in FIG. 2A and FIG. 2B is equivalent to an aircraft flying upside down; since this is known to be an inefficient way to generate lift, it follows that the negatively angled side-fin setting will compromise the basic functions of the side-fin(s), which, as will be appreciated by persons skilled in the art after reading the disclosure which follows below, are as follows: the negative toe-in angle of the side-fins improves directional stability when the rider's weight is evenly balanced on the board; when the rider leans to turn the side-fin functions as a deflected rudder and aids the board's initial rotation and, as the prior art tri-fin (shown in FIG. 2) rotates further in the turn, the angle of the water flow changes so that it is striking the “underside” of the fin(s), which places the fins of the board at a high, “flying” angle of attack and, much like a sail, accelerates the board forward.\nFIG. 3A shows the rotation of the board in more detail: in the diagram depicted, the rider's weight shift when leaning in a turn creates a yawing moment YM that, in relation to the board's original heading H, changes the angle of the “apparent” water-flow F striking the fins, and places the fins at a higher angle of attack. (Note: the term “apparent” water-flow is used in the same manner as the term “apparent wind” is used in sailing—from the board's perspective, the water is “apparently” moving, although the actual angle of the water-flow striking the fins is caused almost entirely by the motion of the board itself.) In the turn shown, an arrow H represents the board's original heading (shown in FIG. 2), while the three arrows running parallel to and in an opposite direction to the first arrow are used to represent the apparent water-flow F resulting from that heading. The board's rotation in the turn is referenced by an imaginary axis of rotation AR, and the arrows at either end of the board depict the direction and rotation R of the nose and tail of the board as the rider, leaning in the turn, shoves the tail in one direction, and causes the nose to move in the opposite direction. The view shows that the movement of the tail as the board rotates causes the fins at the rear of the board to be placed at a higher angle of attack relative to the water-flow F, and increases their potential “lift.” (The “lift” is depicted here as the pressure field described above. In addition, the rotation of the board and angle of attack of the fins may be better visualized if the view is assumed to be from the rider's perspective with the deck or top surface of the board transparent.)\nFIG. 3B depicts a very early, and largely unsuccessful, split-tailed fish style board of the prior art, and shows that the same rotation on a board with a wide tail and correspondingly wide fin-spacing will place the side-fins at a higher angle of attack. The added problem is that on a wide-tailed board the side-fins are further away from the rider's feet—because the rider controls the board through weight shifts that are transmitted through the feet, it follows that a wide side-fin spacing will increase the moment arm of the side-fin, and that the added leverage will lead to control problems since the rider will be less able to counter the reverse yaw of the side-fins and maintain control of the board through a turn.\nFrom the preceding discussion, it will also be apparent that increasing the length of the board 10f or the speed at which it is ridden will exacerbate the problems outlined above. Persons knowledgeable in board design will note that the early double-finned fish style boards, which originally used the parallel side-fin setting shown in FIG. 3B and had large, low aspect ratio keel type fins, were limited to roughly five and a half feet in length. Although these boards at times exhibited exceptional speed in smaller surf, they became difficult or impossible to control at higher speeds in larger, faster waves, where the size of the board was typically increased. As a result, the parallel side-fin setting shown in FIG. 3B was quickly abandoned in favor of the negatively angled side-fin setting of the prior art. The early twin-fin style boards of the same era (not depicted) were also notoriously prone to tracking problems, particularly in larger surf. As will be appreciated by those of skill in the art after reading the disclosure below, this was due to the wide spacing of the side-fins, which were placed near the extreme edge of the very wide square tail and far from the rider's feet.\nTherefore, when comparing the modern prior art tri-fin depicted in FIG. 2 and FIG. 3A to the early, wide-tailed fish of FIG. 3B, it can be seen that the design modifications have comprised a considerable narrowing of the tail; the side-fin placement has moved further forward on the board; and the side-fins are now universally set at a negative angle of attack. These design changes have had the effect of eliminating prior art control problems, but without first identifying their cause—the prior art tri-fin, which is considered to be a fast, exceptionally maneuverable board, retains the inherent drawbacks of the negatively angled side-fin setting, and suffers from seriously compromised performance and considerable unnecessary drag as a result.\nAccordingly, much room remains for improvement in the structure and placement of fins and foils on surfboards."}
{"text": "Ambulatory physiological monitoring has many applications in sports, medicine, industry and the military. In all these applications, it is advantageous to configure sensors on an comfortable, unobtrusive garment to be worn by a monitored individual.\nAmong the known technologies for physiological monitoring is inductive plethysmography (“IP”). IP applied to respiratory monitoring, respiratory IP (“RIP”), has been shown to be a reliable and robust monitoring technology. RIP devices and processing methods are described in U.S. patents and US patent applications including: U.S. Pat. No. 6,047,203 issued Apr. 4, 2000; U.S. Pat. No. 6,551,252 issued Apr. 22, 2002; and U.S. patent application Ser. No. 10/822,260 filed Apr. 9, 2004. These patents are included by reference herein in their entireties for all purposes.\nKnown physiological monitoring garments have been configured in the form of clothing for significant portions of a subject's body, such as shirts, pants, and the like. Known clothing items of such type are not advantageous in certain applications of physiological monitoring, for example, in situations where strenuous activity is possible or when personal garment size fitting is not possible."}
{"text": "In magnetic storage systems for computers, digital data serves to modulate the current in a read/write head coil so that a sequence of corresponding magnetic flux transitions are written onto a magnetic medium in concentric tracks. To read this recorded data, the read/write head passes over the magnetic medium and transduces the magnetic transitions into pulses in an analog signal that alternate in polarity. These pulses are then decoded by read channel circuitry to reproduce the digital data.\nDecoding the pulses into a digital sequence can be performed by a simple peak detector in a conventional analog read channel or, as in more recent designs, by using a discrete time sequence detector in a sampled amplitude read channel. Discrete time sequence detectors are preferred over simple analog pulse detectors because they compensate for intersymbol interference (ISI) and, therefore, are less susceptible to noise. As a result, discrete time sequence detectors increase the capacity and reliability of the storage system. There are several well known discrete time sequence detection methods including discrete time pulse detection (DPD), partial response (PR) with Viterbi detection, maximum likelihood sequence detection (MLSD), decision-feedback equalization (DFE), enhanced decision-feedback equalization (EDFE), and fixed-delay tree-search with decision-feedback (FDTS/DF).\nIn conventional peak detection schemes, threshold crossing or derivative information, implemented in analog circuitry, is normally used to detect peaks in the continuous time analog signal generated by the read head. The analog read signal is \"segmented\" into bit cell periods and interpreted during these segments of time. The presence of a peak during the bit cell period is detected as a \"1\" bit, whereas the absence of a peak is detected as a \"0\" bit. The most common errors in detection occur when the bit cells are not correctly aligned with the analog pulse data. Timing recovery, then, adjusts the bit cell periods so that the peaks occur in the center of the bit cells on average in order to minimize detection errors. Since timing information is derived only when peaks are detected, the input data stream is normally run length limited (RLL) to limit the number of consecutive \"0\" bits.\nDetection errors are also caused by intersymbol interference (ISI). Storage density is directly related to the number of pulses stored on the storage medium, and as the pulses are packed closer together in the effort to increase data density, they eventually interfere with each other resulting in intersymbol interference. This interference can cause a peak to shift out of its bit cell, or its magnitude to decrease, and result in a detection error. The ISI effect is reduced by decreasing the data density or by employing an encoding scheme to ensure that a minimum number of \"0\" bits occur between \"1\" bits. For example, a (d,k) run length limited (RLL) code constrains to d the minimum number of \"0\" bits between \"1\" bits, and to k the maximum number of consecutive \"0\" bits. A typical RLL code is a (1,7) 2/3 rate code which encodes 8 bit data words into 12 bit codewords to satisfy the (1,7) constraint.\nSampled amplitude detection, such as partial response (PR) with Viterbi detection, allows for increased data density by compensating for intersymbol interference. Unlike conventional peak detection systems, sampled amplitude recording detects digital data by interpreting, at discrete time instances, the actual value of the pulse data. The analog pulses are sampled at the baud rate and the digital data is detected from these discrete time sample values. A discrete time sequence detector, such as a Viterbi detector, interprets the discrete time sample values in context to determine a most likely sequence for the data. In this manner, the effect of ISI can be taken into account during the detection process thereby decreasing the probability of a detection error. After processing a consecutive sequence of sample values, the sequence detector compensates for ISI by selecting the most likely digital sequence associated with the sample values. This increases the signal to noise ratio and, for a given (d,k) constraint, allows for significantly higher data density as compared to conventional analog peak detection read channels.\nThe application of sampled amplitude techniques to digital communication channels is well documented. See Y. Kabal and So Pasupathy, \"Partial Response Signaling\", IEEE Trans. Commun. Tech., Vol. COM-23, pp.921-934, September 1975; and Edward A. Lee and David G. Messerschmitt, \"Digital Communication\", Kluwer Academic Publishers, Boston, 1990; and G. D. Forney, Jr., \"The Viterbi Algorithm\", Proc. IEEE, Vol. 61, pp. 268-278, March 1973.\nApplying sampled amplitude techniques to magnetic storage systems is also well documented. See Roy D. Cideciyan, Francois Dolivo, Walter Hirt, and Wolfgang Schott, \"A PRML System for Digital Magnetic Recording\", IEEE Journal on Selected Areas in Communications, Vol. 10 No. 1, January 1992, pp.38-56; and Wood et al, \"Viterbi Detection of Class IV Partial Response on a Magnetic Recording Channel\", IEEE Trans. Commun., Vol. Com-34, No. 5, pp. 454-461, May 1986; and Coker Et al, \"Implementation of PRML in a Rigid Disk Drive\", IEEE Trans. on Magnetics, Vol. 27, No. 6, November 1991; and Carley et al, \"Adaptive Continous-Time Equalization Followed By FDTS/DF Sequence-Detection\", Digest of The Magnetic Recording Conference, Aug. 15-17, 1994, pp. C3; and Moon et al, \"Constrained-Complexity Equalizer Design for Fixed Delay Tree Search with Decision Feedback\", IEEE Trans. on Magnetics, Vol. 30, No. 5, September 1994; and Abbott et al, \"Timing Recovery For Adaptive Decision Feedback Equalization of The Magnetic Storage Channel\", Globecorn'90 IEEE Global Telecommunications Conference 1990, San Diego, Calif., November 1990, pp.1794-1799; and Abbott et al, \"Performance of Digital Magnetic Recording with Equalization and Offtrack Interference\", IEEE Transactions on Magnetics, Vol. 27, No. 1, January 1991; and Cioffi et al, \"Adaptive Equalization in Magnetic-Disk Storage Channels\", IEEE Communication Magazine, February 1990; and Roger Wood, \"Enhanced Decision Feedback Equalization\", Intermag'90."}
{"text": "Light diffusing films are currently widely used in many optical, lighting, and display applications. Many of the most efficient diffusers utilize engineered microstructures to effectively distribute the light in a desired manner. Most of these diffusing films are formed by either molding or casting a polymer against a mold surface having relief structures present to produce the appropriate structure on the film surface. These structures are usually made of non-elastomeric polymeric materials, such as polyacrylates, polystyrenes, polycarbonates and polyolefins.\nA problem arises when the required diffusing microstructures are of a shape that it becomes fragile and thereby easily damaged by routine handling of the film in assembly processes. Minor surface rubs can result in permanent, visible scratches and abrasions on the diffusing film rendering it useless for many applications or at least cosmetically. Therefore, a need exists for an improved light diffusing film."}
{"text": "Generic lighting control consoles are used for controlling lighting systems such as are used for example in theaters and/or on concert stages. These lighting systems normally include a large number of lighting devices, such as stage spotlights, and in their turn the lighting devices themselves can be switched between a wide variety of lighting states, for example different colors. The various lighting devices with their different lighting states are controlled by programmed parameters in the lighting software of the lighting control console."}
{"text": "The field of art to which this invention pertains is the treatment of aqueous streams containing water-soluble inorganic sulfide compounds. More specifically, the invention is directed toward a method for treating an aqueous stream containing a water-soluble, inorganic sulfide compound."}
{"text": "In a conventional technology related to an electronic control device equipped in vehicles that can control in-vehicle equipment, module type electronic equipment disclosed in JP-A-8-274440 is known.\nThis module type electronic equipment is constituted so that a power source module can be inserted to and detached from a front side of a card cage or a power supply module part.\nA power supply module connector is automatically connected with a power supply connector in a backplane when the power source module is completely inserted into a card receptor, and the connection is automatically disconnected when the power source module is detached from the card cage.\nMoreover, a daughter circuit card is constituted similar to the connection structure as the above, and an operation power is supplied easily without wiring each other with the power supply relative to the power source module and the daughter circuit card.\nAn electronic system disclosed in JP-A-2004-246615 is a device adopted for personal computers or servers, and is constituted with a plurality of electronic modules, circuit board modules for control, and power source modules, etc. installed in a case.\nA CPU that is a heat generating element and other heat generating elements are installed on each surfaces of the installed electronic modules together with a hard disk etc., and each electronic module is a part of the computer.\nMoreover, a respective necessary electrical connection between each electronic module and the circuit board module or the power source module, etc. of the device side is performed by inserting (press-fitting) the connector into a connector of a common wiring circuit board.\nA liquid circulation loop is constituted by inserting (press-fitting) liquid connectors of each electronic module into liquid connectors of a liquid cooling jacket at the time of the electric connection, and cooling of each electronic module becomes possible by a coolant that flows in the liquid circulation loop.\nBy the way, in recent years, the number of electronic control unit (ECU) required for in-vehicle equipment such as an engine and a transmission, etc. equipped in a vehicle increases in order to cope with the demand of making the vehicle with high performance.\nHowever, when the number of the in-vehicle ECU increases, it is necessary to provide an assembly space and a cooling space for each ECU, thus causes a problem that required capacity for equipment in the vehicle increases.\nMoreover, the number of wire harnesses that electrically connect between each ECU and between the ECU and power supplies increases according to the increase of ECUs, thus causes a problem that required capacity for equipment in the vehicle further increases.\nSince the current used is low and the function of the electronic module is limited, a composition adopted in a personal computer or a server as shown in the above-mentioned JP-A-2004-246615 is insufficient to deal with the electronic control unit for the vehicle that has various functions increasing.\nMoreover, the composition of the above-mentioned JP-A-8-274440 has a similar problem."}
{"text": "This application claims priority from Japanese Patent Application No. 2002-084407 filed Mar. 25, 2002, which is incorporated hereinto by reference.\n1. Field of the Invention\nThe present invention relates to an ink jet print head and an ink jet printing apparatus for ejecting ink onto a print medium to form an image thereon.\n2. Description of the Related Art\nAn ink jet printing apparatus forms an image on a print medium by ejecting ink droplets from a print head mounted in an apparatus body onto the print medium, with the ink droplets adhering to the print medium and fixing in it to produce their intended colors. Recent years have seen a proliferation of a so-called serial scan type ink jet printing apparatus. In this type, an image is formed by alternately repeating two operationsxe2x80x94a printing scan for scanning the print head over the print medium to eject ink onto the medium and a paper feed for moving the print medium or the print head relative to each other in a direction perpendicular to a printing scan direction. The serial scan type ink jet printing apparatus, however, has the following drawback.\nIn the serial scan type apparatus, a single printing scan can only produce an image of a predetermined printing width for at least one color of ink (this single printing operation is referred to also as a xe2x80x9cone-pass printingxe2x80x9d). Hence, to form an image over the entire print medium requires performing a plurality of printing scans. When in such a system an image of high duty is to be formed, a problem may occur that a boundary portion between an image area formed on the print medium in a certain printing scan and an adjoining image area formed in another printing scan appears light in density.\nThis problem is considered to occur in the following mechanism. FIGS. 7A to 7D are schematic views showing how an image of high duty is formed during the one-pass printing, as seen in the print head scanning direction. In the figure, reference number 1 represents a print head, 2 a print medium, and e a column of nozzles (also referred to as a xe2x80x9ccolumn of ejection openingsxe2x80x9d) for ejecting ink droplets.\nFIG. 7A shows ink droplets adhering to a print medium which were ejected in one printing scan. In the figure, p1 denotes an ink adhering to the print medium. With the elapse of time the ink on the print medium soaks into the medium and fixes there. FIG. 7B illustrates this state and p2 denotes the ink that has soaked into the print medium and fixed there. After the printing scan, the print medium is fed in a direction perpendicular to the printing scan direction of the print head (this operation is called a line feed) and the next printing scan is performed. FIG. 7C shows a state in which the line feed and the second printing scan have been performed. In the figure, a distance that the print medium was fed is indicated by an arrow. This line feed distance is equal to the length of the nozzle column of the print head. Ink droplets adhering to the print medium that were ejected in the second printing scan are indicated by p3.\nFIG. 7D shows a state in which the ink that landed on the print medium during the second printing scan has soaked and fixed with elapse of time. As shown by p1 of FIG. 7A and p3 of FIG. 7C, the ink that has just landed on the print medium and has not yet soaked into and fixed in the print medium forms an ink surface that is low at ends and bulges at a center. This is a common phenomenon produced by a surface tension of the ink. In this state, the ink penetrates and fixes in the print medium. Therefore, as indicated by p2 and p4 of FIGS. 7B and 7D, in an image area formed by each printing scan, an amount of ink that fixes at the end portions is less than at other portions and the color of that portions tends to be lighter. Thus, when the printing scan is repeated a plurality of times to form an image of high duty, the end portions of an image area formed by each printing scan appear light. That is, the boundary portions between adjoining image areas are printed lighter than other portions, giving rise to a problem of light stripes showing up in the printed image. In the case of a black ink in particular, since its penetration capability is generally low, it tends to produce a greater density difference between a dot center and a dot end than do color inks. This may result in boundary portions between adjoining image areas printed by different printing scans appearing lighter and in the worst case showing up as white horizontal stripes.\nA possible measure to deal with this problem may involve making the line feed distance shorter than the printing width or nozzle column length of the print head. One such example is to design a nozzle-to-nozzle interval (nozzles may also be referred to as xe2x80x9cejection openingsxe2x80x9d) somewhat longer than normal. As a result, the length of the nozzle column used for the one-pass printing becomes somewhat longer than the line feed distance, producing the following advantages.\nFIGS. 8A to 8D show dots ejected from a nozzle column with a longer-than-normal nozzle-to-nozzle interval. In the figure, reference numeral 1 represents a print head, 2 a print medium and e a nozzle column for ejecting ink droplets. FIG. 8A shows ink droplets adhering to the print medium which were ejected in one printing scan. In the figure, p1 denotes an ink adhering to the print medium. With the elapse of time the ink on the print medium soaks into the print medium and fixes there. FIG. 8B illustrates this state and p2 denotes the ink that has soaked into the print medium and fixed there. After the first printing scan, the print medium is fed (line feed) in a direction perpendicular to the printing scan direction of the print head and the next printing scan is performed. Because the nozzle-to-nozzle interval of the print head is set somewhat longer than normal, the line feed distance is shorter than the nozzle column length e. FIG. 8C shows a state in which the line feed and the second printing scan have been performed. In the figure, a distance that the print medium was fed is indicated by an arrow and, as described above, is somewhat shorter than the nozzle column length e of the print head. Ink droplets adhering to the print medium that were ejected in the second printing scan are indicated by p3. Then, the ink that landed on the print medium during the second printing scan also sinks and fixes in the print medium over time, as shown in FIG. 8D.\nSince the line feed distance shown in FIG. 7C is equal to the nozzle column length or a difference between the line feed distance and the nozzle column length is smaller than that of FIG. 8C, a comparison between FIG. 8D and FIG. 7D shows that an overlap between p2 and p4 is somewhat larger in FIG. 8D than in FIG. 7D. Thus, as shown in FIG. 8D, the above-described problem that a boundary portion between an image area formed on the print medium by a printing scan and an adjoining image area formed by another printing scan appears lighter than other portions is less likely to occur.\nPrinting apparatus capable of printing color inks as well as black ink are available in recent years. Some of these printing apparatus have a black ink nozzle column set longer than other color ink nozzle columns in order to reduce a time taken by the printing operation using only the black ink as in a document printing. In this arrangement, when printing is done using only the black ink, all the nozzles of the black ink nozzle column are used, whereas during color printing, only that part of the black ink nozzle column which is almost equal in length to other color ink nozzle columns is used. In such a printing apparatus, in which the length of that nozzle portion in the entire nozzle column which is used for printing is changed according to an image being formed, a problem may arise that lighter horizontal stripes will show up in a printed image at boundaries between adjoining image areas formed on a print medium by separate printing scans, depending on the length of the nozzle portion used for printing. This problem will be explained as follows.\nReferring to FIG. 2 and FIG. 3, reference number 1 denotes a print head, 3 a nozzle column for ejecting a black ink, and 4 nozzle columns for ejecting color inks. To solve the problem described above, the black ink nozzle column is formed longer than the color ink nozzle columns. In the black nozzle column 3, the entire nozzles are represented as a nozzle portion e and a part of the nozzle column is denoted a nozzle portion b. The nozzle portion b has one-half the length of the nozzle portion e. The entire nozzles arrayed in each of the color ink nozzle columns are represented as a nozzle portion a. The number of nozzles in the nozzle portion a counted in the column direction is equal to that of the nozzle portion b.\nFIG. 2 is a schematic view showing an operation of the printing apparatus when an image is formed using only a black ink. When an image is formed using only the black ink, the whole black nozzle column (nozzle portion e) is used as described above. In the figure, (f1)-p1 represents a position relative to the print head of an image formed with the black ink in one printing scan. This is followed by a line feed of a predetermined distance in a direction indicated by LF. The line feed distance is shorter than the length of the nozzle portion e. The printed image p1 moves to a position (f2)-p1. After this, another printing scan is performed to form an image (f2)-p2.\nFIG. 3 is a schematic view showing an operation of the printing apparatus when an image is formed using a black ink and color inks. As described above, when an image is formed using color inks as well as a black ink, the nozzle portion b of the black nozzle column and the nozzle portion a of the color nozzle columns are used. In the figure, (f1)-p1 represents a position relative to the print head of an image formed with the black ink in one printing scan. After this, a line feed of a predetermined distance is carried out in the direction of LF, moving the printed image p1 to a position (f2)-p1. This is followed by another printing scan to form an image at a position (f2)-p1 using color inks and an image at a position (f2)-p2 using a black ink. As a result, in the (f2)-p1 area the image forming using the black ink and the color inks is completed.\nWhether an image is to be made using only a black ink or both a black ink and color inks is determined based on image data sent from a host computer. A printer driver running on the host computer displays an operation window for the user to select either a color printing or a black-only printing. When the user makes a selection on the operation window, the printer driver sends a color printing instruction or a black-only printing instruction along with image data to the printing apparatus. The printing apparatus determines the operations of various driving units according to the instruction received. Another arrangement is also available in which, rather than the user selecting either a color printing or a black-only printing, the printing apparatus checks the image data transferred from the host to make a decision. Still another arrangement is available in which a detailed control is performed to switch the black nozzle operation between a long nozzle portion and a short nozzle portion of the black nozzle column according to the image data in each page. That is, in an area of each page to be printed with only a black ink a long black nozzle portion, i.e., entire black nozzle column, is used and, in an area to be printed with color inks as well, a short black nozzle portion equal in length to the color nozzle columns is used.\nIn a printing apparatus with a means to change the length of a black nozzle portion to be used for printing, it has been proposed to set a nozzle interval a predetermined amount longer than normal to deal with the aforementioned problem of light density portions showing up in a printed image at boundaries between image areas printed by separate printing scans. As explained earlier in conjunction with FIG. 2 and FIG. 3, the length of an activated portion of the black nozzle column differs between the black-only printing and the color printing. Therefore, the difference between the line feed distance and the width (in the line feed direction) of a black printed area also varies. More specifically, the black nozzle column is set somewhat longer than normal by expanding the nozzle intervals uniformly. If it is assumed that the black nozzle column is set longer by t than the normal nozzle column length s, an entire length of the nozzle column is s+t. In a black-only printing, the entire black nozzle column is used and, if the line feed distance is assumed to be s, image areas printed by separate printing scans overlap each other over a distance of t. In a color printing, only the nozzle portion b of the black nozzle column is used, that is, only one-half of the black nozzle column is used. Then, the length of the nozzle portion b is 1/2xc2x7(s+t). Suppose that the line feed distance is s/2. The difference between the line feed distance and the length of the nozzle portion b is only t/2. Thus the overlap between the image areas is only t/2. This means that, if the nozzle interval is expanded to ensure an enough overlap during the black-only printing, the color printing cannot secure a sufficient overlap. Conversely, if the nozzle interval is set so as to cause a sufficient overlapping during the color printing, the amount of overlap at the boundary portions between separate printing scans becomes too large, giving rise to a problem that the overlapped portions may look darker than other portions.\nIn light of the conventional problems described above, it is an object of the present invention to provide an ink jet print head and an ink jet printing apparatus which can produce a good printed result at all times at boundary portions between image areas printed by separate printing scans even in ink jet printing apparatus in which a range of use of the nozzle column varies according to the printing condition.\nIn one aspect, the present invention provides an ink jet print head having a plurality of nozzles arrayed in a predetermined direction to form a nozzle column, wherein the nozzle column ejects ink droplets, the ink jet print head comprising: a long nozzle column portion formed in a predetermined portion of the nozzle column, the long nozzle column portion having a wider nozzle interval than those in other portions of the nozzle column.\nIn another aspect, the present invention provides an ink jet printing apparatus comprising: a ink jet print head having a plurality of nozzles arrayed in a predetermined direction to form a nozzle column, the nozzle column being adapted to eject ink droplets; wherein the ink jet print head is scanned over a print medium a plurality of times in a direction different from the direction of array and a printing scan and a line feed are performed to print on a predetermined image area on the print medium, the printing scan ejecting ink droplets onto the print medium during each scan and the line feed feeding, between each of the plurality of scans, the print medium and the ink jet print head relative to each other in a direction different from the scan direction of the ink jet print head; wherein a portion of the nozzle column in the ink jet print head is a long nozzle column portion whose nozzle-to-nozzle interval is wider than that in another portion of the nozzle column; wherein a width in the line feed direction of each image area printed by a single printing scan of the ink jet print head is longer than a distance that the print medium is fed by one print feed.\nWith this construction, by arranging the nozzles in the nozzle column such that, in only that portion of the nozzle column always used in any printing condition, such as color printing and black-only printing, its nozzles have a wider nozzle-to-nozzle interval than those of other nozzle portions, the width of each image area printed by a single printing scan can be made a predetermined amount longer than the line feed distance at all times. This arrangement can produce a printed result in which adjoining image areas printed by separate printing scans overlap each other at their boundary portions by a predetermined amount in whatever printing condition.\nThe above and other objects, effects, features and advantages of the present invention will become more apparent from the following description of embodiments thereof taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a technique of lithography and, particularly to a manufacturing method for semiconductor device wherein a stepwise taper is formed in a contract hole, that is, a opening portion, which is formed in an insulating film and metal, by the combination of a resist patterning process and a dry etching process.\n2. Description of the Related Art\nAs is known, the degree of integration of the semiconductor element has been improved with the advance of a fine process technique of the semiconductor element. For the purpose of fining, the so-called multilayered wiring element has been generally used. In order to control leak current generated in the element, the necessary an interlayered insulating film is formed to have the thickness of 1 .mu.m or more regardless of quality. Thereby, the depth of the opening portion, which is formed in the interlayered insulating film, is increased.\nAs one of the fine process techniques, the frequency in use of a dry etching such as a reactive ion etching, in which an anisotropic etching is performed, is increased. For this reason, the side wall of the opening portion, which is formed in the interlayered insulating film, becomes almost vertical. Therefore, in a case where the inside of the opening portion cannot be completely embedded with the metal layer for wiring, which is deposited on the interlayered insulating film, the connection between the electrode in the opening portion and the external wiring layer is defective. To overcome this problem, there is known a method for forming taper in the opening portion contacting the wiring metal layer, that is, 20 to 30% of the entire depth, by using the combination of a chemical reaction or a physical etching.\nAlso, an image reversal process which can be used as the fine process technique was disclosed in Vacuum, Vol. 29, No. 12, 1989, page 585. Specifically, a mask, in which a layer interrupting light is formed, is superimposed at a predetermined position of a positive resist of a novolak type, an exposing process, which is a part of a photo engraving process, is provided therein, so that a photosensitive portion is formed therein other than the light interrupting layer of the mask. However, if this lamination layer is processed in ammonia atmosphere, the exposed and photosensitive indene carboxylic acid generates decarboxylic acid, insoluble indene is formed in alkali developer. If this phenomenon is used, a pattern to be originally formed is reversed. FIG. 1 illustrates an example of forming taper by the combination of the chemical etching and the physical etching. More specifically, a wiring layer forming of Al or Al alloy or an electrode 2 is directly or indirectly deposited on a semiconductor substrate 1 via a coating layer, and an interlayered insulating layer 3 is entirely deposited thereon.\nThen, an opening portion 4 is formed on the interlayered insulating layer 3 by CDE process as a chemical etching, in which radical, which is generated by plasma generated by a magnetron tube, is moved to the distant position and RIE process as a physical etching. The edge of the opening portion 4 is etched in a certain horizontal and vertical isotropic direction by CDE process. Due to this, as is obvious from the drawing, a taper 6 is formed and a double edge is formed in positioning marks 5, which are used in an aligner process, as shown in FIG. 2. FIG. 2 is a plane view of the opening portion 4 for the positioning marks 5 seen from the above.\nMoreover, FIG. 3 is a cross sectional view showing the opening portion 4 formed by providing image reversal process in the interlayered insulating layer 3. However, the side wall of the opening portion 4 is not vertically formed, and the double edge is formed in the positioning marks which are formed at the same with the above process. This causes the deterioration of the accuracy of the positioning.\nIn a case where the taper is formed in the opening portion 4 by the combination of the chemical etching and the physical etching, the positioning marks 5 are shaped similarly and the edges of the positioning marks look double as shown in FIG. 2. In general, to detect the positioning marks 5, reflected light of the irradiated laser is electrically compared between a setting coordinate and a measuring coordinate. Then, if the edges are discriminated several times, a plurality of detection signals are obtained, so that the positioning step fails or the accuracy of positioning is deteriorated.\nMoreover, if the opening portion and the positioning marks are formed by use of the image reversal process, both the opening portion and the positioning marks are simultaneously formed, so that the taper 6 is formed in the positioning marks. As a result, since a detection signal having a clear waveform is not obtained, the accuracy of the positioning is deteriorated. In addition, regarding etching, it is difficult to set conditions such as an angle of a reverse taper and controllability and there is problem in reproductivity."}
{"text": "1. Technical Field\nThe present invention is directed to an improved data processing system. More specifically, the present invention is directed to an apparatus and method for implementing multicast on a system area network channel adapter associated with a logically partitioned (LPAR) data processing system, with no visibility to either the Fabric Manager (Subnet Manager) or other Fabric Participants, that LPAR techniques are being employed.\n2. Description of Related Art\nInfiniBand (IB), which is a form of System Area Network (SAN), defines a multicast facility that allows a Channel Adapter (CA) to send a packet to a single address and have it delivered to multiple ports. Each multicast group is assigned a unique address, and end-nodes that wish to participate in a multicast group do so via a ‘Join’ process initiated by the candidate participant with the Subnet Manager. The InfiniBand architecture is described in the InfiniBand standard, which is available at http://www.infinibandta.org and also hereby incorporated by reference.\nWith the InfiniBand architecture, the CA sending the multicast packet may be a Host Channel Adapter (HCA) or a Target Channel Adapter (TCA). A multicast packet is sent to all ports of a collection of ports called a multicast group. These ports may be on the same or different nodes in the SAN. Each multicast group is identified by a unique Local Identifier (LID) and Global Identifier (GID). The LID is an address assigned to a port which is unique within the subnet. The LID is used for directing packets within the subnet. The GID is a 128-bit identifier used to uniquely identify a port on a channel adapter, a port on a router, or a multicast group, across all infiniband subnets. The LID and GID are in the Local Route Header (LRH) and Global Route Header (GRH), respectively, of the IB packet. The LRH is present in all IB packets and is an address used for routing IB packets through switches within a subnet. The GRH is present in IB packets which are either multicast packets, or which are targeted to destinations outside the originator's local subnet and is used as an address for routing the packets when the packets traverse multiple subnets.\nAn IB management action via a Subnet Management Packet (SMP) is used when a node joins a multicast group, and at that time the LID of the port on the node is linked to the multicast group. A subnet manager then stores this information in the switches of the SAN using SMPs. The subnet manager via SMPs tells the switches the routing information for the various multicast groups, and the switches store that information, so that the switches can route the multicast packets to the correct nodes. When a node is going to send a packet to the multicast group, it uses the multicast LID and GID of the group to which it wants the packet to be delivered. The switches in the subnet detect the multicast LID in the packet's Destination LID (DLID) field and replicates the packet, sending it to the appropriate ports, as previously set up by the subnet manager.\nIt is the Subnet Manager's job to look at the topology and adjust the multicast forwarding tables of each applicable switch in the fabric such that a member of a multicast group will receive a multicast packet sent to that Multicast Group address.\nWithin a CA, one or more Queue Pairs (QPs) may be registered to receive a given multicast address. IB allows for the number of QPs within a CA that can be registered for the same address to be only limited by the particular implementation. The registration process is done via the IB verb interface. The verb interface is an abstract description of the functionality of a Host Channel Adapter. An operating system exposes some or all of the verb functionality through its programming interface.\nWhen the CA recognizes a multicast packet, the CA must somehow distribute the packet to all the registered QPs within that CA. This must be done in an efficient manner. How this is done is not specified by the InfiniBand Architecture (IBA).\nCommonly-owned co-pending Published U.S. Patent Application No. 2003/003426 of Beukema et al., application Ser. No. 09/925,578, filed Aug. 9, 2001, which is incorporated herein by reference, describes a system for implementing multicast on an Infiniband CA. However, the solution described in the Beukema application does not address the additional complexity associated with a logically-partitioned (LPAR) data processing system.\nWhen implementing LPAR, it is advantageous that each Operating System believes that it has control of a single CA. This is further substantiated by the requirement to maintain transparency to the Subnet Manager and other end-nodes, i.e., neither of these must operate any differently when talking to an LPAR end-node vs. a non-LPAR end-node. In order to achieve this, each LPAR sees a logical CA. The ports on this logical CA are assigned LIDs, just like real ports. In addition, packets coming into the ‘real’ port of a CA effectively see a logical switch. This logical switch has a set of logical Multicast Forwarding Tables that the Subnet manager will set up.\nIn an LPAR computing environment, a single data processing system is “virtualized” to multiple software partitions, each representing a different instance of an operating system. An LPAR data processing system thus functions as if it were several separate machines, though the “machines” (generally unbeknownst to each other) share a common hardware platform. LPAR systems are well suited for situations in which multiple computing platforms are needed, but the additional expense and inconvenience of installing and maintaining multiple physical hardware platforms is undesirable. In particular, it would be beneficial if a CA for a SAN such as Infiniband could be shared among multiple partitions of an LPAR system."}
{"text": "Medical ceiling pendant is the essential medical device in the modern operating room in the hospital, which is mainly used for terminal-transferring the medial gas, such as oxygen supply, compressed air, nitrogen, etc., in the operating room and the intensive care unit, the lower part of the medical ceiling pendant is provided with a load platform, which facilitates placing some surgical instruments or equipments and meters. As the huge structure size of the medical ceiling pendant, the arms and the box body are separately packaged during the prior packaging and transportation, such packaging process has the defects as follows: 1. there needs two packing box for one medical ceiling pendant, thus resulting in more consumables and higher logistical cost, 2. the gas tubes and the wires on the box body are required to pass through the ceiling pendant arm at the customer's field, thus resulting in time-consuming and labor-consuming, 3. it requires some skills to pass through the gas tubes and the wires in field, thus unable to guarantee the quality."}
{"text": "An important large family of enzymes is the protein kinase enzyme family. Currently, there are about 500 different known protein kinases. However, because three to four percent of the human genome is a code for the formation of protein kinases, there may be many thousands of distinct and separate kinases in the human body. Protein kinases serve to catalyze the phosphorylation of an amino acid side chain in various proteins by the transfer of the γ-phosphate of the ATP-Mg2+ complex to said amino acid side chain. These enzymes control the majority of the signaling processes inside cells, thereby governing cell function, growth, differentiation and destruction (apoptosis) through reversible phosphorylation of the hydroxyl groups of serine, threonine and tyrosine residues in proteins. Studies have shown that protein kinases are key regulators of many cell functions, including signal transduction, transcriptional regulation, cell motility, and cell division. Several oncogenes have also been shown to encode protein kinases, suggesting that kinases play a role in oncogenesis. These processes are highly regulated, often by complex intermeshed pathways where each kinase will itself be regulated by one or more kinases. Consequently, aberrant or inappropriate protein kinase activity can contribute to the rise of disease states associated with such aberrant kinase activity. Due to their physiological relevance, variety and ubiquitousness, protein kinases have become one of the most important and widely studied family of enzymes in biochemical and medical research.\nThe protein kinase family of enzymes is typically classified into two main subfamilies: Protein Tyrosine Kinases and Protein Serine/Threonine Kinases, based on the amino acid residue they phosphorylate. The serine/threonine kinases (PSTK), includes cyclic AMP- and cyclic GMP-dependent protein kinases, calcium and phospholipid dependent protein kinase, calcium- and calmodulin-dependent protein kinases, casein kinases, cell division cycle protein kinases and others. These kinases are usually cytoplasmic or associated with the particulate fractions of cells, possibly by anchoring proteins. Aberrant protein serine/threonine kinase activity has been implicated or is suspected in a number of pathologies such as rheumatoid arthritis, psoriasis, septic shock, bone loss, many cancers and other proliferative diseases. Accordingly, serine/threonine kinases and the signal transduction pathways which they are part of are important targets for drug design. The tyrosine kinases phosphorylate tyrosine residues. Tyrosine kinases play an equally important role in cell regulation. These kinases include several receptors for molecules such as growth factors and hormones, including epidermal growth factor receptor, insulin receptor, platelet derived growth factor receptor and others. Studies have indicated that many tyrosine kinases are transmembrane proteins with their receptor domains located on the outside of the cell and their kinase domains on the inside. Much work is also under progress to identify modulators of tyrosine kinases as well.\nNuclear factor κB (NF-κB) belongs to a family of closely related dimeric transcription factor complexes composed of various combinations of the Rel/NF-κB family of polypeptides. The family consists of five individual gene products in mammals, RelA (p65), NF-κB1 (p50/p105), NF-κB2 (p49/p100), c-Rel, and RelB, all of which can form hetero- or homodimers. These proteins share a highly homologous 300 amino acid “Rel homology domain” which contains the DNA binding and dimerization domains. At the extreme C-terminus of the Rel homology domain is a nuclear translocation sequence important in the transport of NF-κB from the cytoplasm to the nucleus. In addition, p65 and cRel possess potent transactivation domains at their C-terminal ends.\nThe activity of NF-κB is regulated by its interaction with a member of the inhibitor IκB family of proteins. This interaction effectively blocks the nuclear localization sequence on the NF-κB proteins, thus preventing migration of the dimer to the nucleus. A wide variety of stimuli activate NF-κB through what are likely to be multiple signal transduction pathways. Included are bacterial products (LPS), some viruses (HIV-1, HTLV-1), inflammatory cytokines (TNFα, IL-1), environmental and oxidative stress and DNA damaging agents. Apparently common to all stimuli however, is the phosphorylation and subsequent degradation of IκB. IκB is phosphorylated on two N-terminal serines by the recently identified IκB kinases (IKK-α and IKK-β). IKK-β is also known as IKK2. Site-directed mutagenesis studies indicate that these phosphorylations are critical for the subsequent activation of NF-κB in that once phosphorylated the protein is flagged for degradation via the ubiquitin-proteasome pathway. Free from IκB, the active NF-κB complexes are able to translocate to the nucleus where they bind in a selective manner to preferred gene-specific enhancer sequences. Included in the genes regulated by NF-κB are a number of cytokines and chemokines, cell adhesion molecules, acute phase proteins, immunoregulatory proteins, eicosanoid metabolizing enzymes and anti-apoptotic genes.\nIt is well-known that NF-κB plays a key role in the regulated expression of a large number of pro-inflammatory mediators including cytokines such as TNF, IL-1β, IL-6 and IL-8, cell adhesion molecules, such as ICAM and VCAM, and inducible nitric oxide synthase (iNOS). Such mediators are known to play a role in the recruitment of leukocytes at sites of inflammation and in the case of iNOS, may lead to organ destruction in some inflammatory and autoimmune diseases.\nThe importance of NF-κB in inflammatory disorders is further strengthened by studies of airway inflammation including asthma, in which NF-κB has been shown to be activated. This activation may underlie the increased cytokine production and leukocyte infiltration characteristic of these disorders. In addition, inhaled steroids are known to reduce airway hyperresponsiveness and suppress the inflammatory response in asthmatic airways. In light of the recent findings with regard to glucocorticoid inhibition of NF-κB, one may speculate that these effects are mediated through an inhibition of NF-κB.\nFurther evidence for a role of NF-κB in inflammatory disorders comes from studies of rheumatoid synovium. Although NF-κB is normally present as an inactive cytoplasmic complex, recent immunohistochemical studies have indicated that NF-κB is present in the nuclei, and hence active, in the cells comprising rheumatoid synovium. Furthermore, NF-κB has been shown to be activated in human synovial cells in response to stimulation with TNF-α or IL-1β. Such a distribution may be the underlying mechanism for the increased cytokine and eicosanoid production characteristic of this tissue. See Roshak, A. K., et al., J. Biol. Chem., 271, 31496-31501 (1996). Expression of IKK-β has been shown in synoviocytes of rheumatoid arthritis patients and gene transfer studies have demonstrated the central role of IKK-β in stimulated inflammatory mediator production in these cells. See Aupperele et al. J. Immunology 1999. 163:427-433 and Aupperle et al. J. Immunology 2001; 166:2705-11. More recently, the intra-articular administration of a wild type IKK-β adenoviral construct was shown to cause paw swelling while intra-articular administration of dominant-negative IKKβ inhibited adjuvant-induced arthritis in rat. See Tak et al. Arthritis and Rheumatism 2001, 44:1897-1907.\nThe NF-κB/Rel and IκB proteins are also likely to play a key role in neoplastic transformation and metastasis. Family members are associated with cell transformation in vitro and in vivo as a result of over expression, gene amplification, gene rearrangements or translocations. In addition, rearrangement and/or amplification of the genes encoding these proteins are seen in 20-25% of certain human lymphoid tumors. Further, NF-κB is activated by oncogenic ras, the most common defect in human tumors and blockade of NF-κB activation inhibits ras mediated cell transformation. In addition, a role for NF-κB in the regulation of apoptosis has been reported strengthening the role of this transcription factor in the regulation of tumor cell proliferation. TNF, ionizing radiation and DNA damaging agents have all been shown to activate NF-κB which in turn leads to the upregulated expression of several anti-apoptotic proteins. Conversely, inhibition of NF-κB has been shown to enhance apoptotic-killing by these agents in several tumor cell types. As this likely represents a major mechanism of tumor cell resistance to chemotherapy, inhibitors of NF-κB activation may be useful chemotherapeutic agents as either single agents or adjunct therapy. Recent reports have implicated NF-κB as an inhibitor of skeletal cell differentiation as well as a regulator of cytokine-induced muscle wasting (Guttridge et al. Science; 2000; 289: 2363-2365.) further supporting the potential of NFκB inhibitors as novel cancer therapies.\nSeveral NF-κB inhibitors are described in C. Wahl, et al. J. Clin. Invest. 101(5), 1163-1174 (1998), R. W. Sullivan, et al. J. Med. Chem. 41, 413-419 (1998), J. W. Pierce, et al. J. Biol. Chem. 272, 21096-21103 (1997).\nThe marine natural product hymenialdisine is known to inhibit NF-κB. Roshak, A., et al., JPET, 283, 955-961 (1997). Breton, J. J and Chabot-Fletcher, M. C., JPET, 282, 459-466 (1997).\nAdditionally, patent applications have been filed on aminothiophene inhibitors of the IKK2, see Callahan, et al., WO 2002030353; Baxter, et al., WO 2001058890, Faull, et al., WO 2003010158; Griffiths, et al., WO2003010163; Fancelli, et al., WO 200198290; Granetto, et al., WO 2003037886; imidazole inhibitors of IKK2, see Callahan, et al., WO 200230423; anilinophenylpyrimidine inhibitors of IKK2, see Kois, et al., WO 2002046171; β-carboline inhibitors of IKK2, see Ritzeler, et al, WO 2001068648, Ritzeler, et al, EP 1134221; Nielsch, et al. DE 19807993; Ritzeler, et al., EP 1209158; indole inhibitors of IKK2, see Ritzeler, et al., WO 2001030774; benzimidazole inhibitors of the IKK2, see Ritzeler, et al., DE 19928424; Ritzeler et al, WO 2001000610; Ritzeler, et al., WO 2004022553; aminopyridine inhibitors of IKK2, see Lowinger, et al, WO 2002024679; Murata, et al, WO 2002024693; Murata, et al., WO 2002044153; aminopyrimidine inhibitors of IKK2, see Bollbuck, et al., WO 2004089913; pyrazole inhibitors of IKK2, see Bergmanis, et al., WO 2003024935; Metz, et al., WO 2003024936; Geng et al., WO 2003027075; Stealey, et al., WO 2003035625; Xu, et al., WO 200307076; Lennon, et al., WO 2003095430; pyrazinone inhibitors of IKK2, see Boys, et al., WO 2005035527; pyrazolaquinazoline inhibitors of IKK2, see Beaulieu, et al., WO 2002028860; Burke et al, WO 2002060386; Burke, et al. US 20030022898; thiophene tricyclic inhibitors of IKK2, see Belema, et al., WO 2003084959; pyrazolopurine inhibitors of IKK2, see Qiu, et al., WO 2004075846; oxazolo and thiazolo pyridine inhibitors of IKK2, see Pitts, et al., WO 2004106293; quinoline inhibitors of IKK2, Browner, et al., WO2002041843, Browner, et al., US 20020161004 and pyridylcyanoguanidine inhibitors of IKK2, see Bjorkling, et al., WO 2002094813, Binderup et al, WO 2002094322 and Madsen, et al., WO 200294265; thienopyridine inhibitors of IKK2, see Cywin, et al., WO 2003103661; Liu, et al., WO 2005035537; benzothiophene inhibitors of IKK2, see Chen et al., WO 2005012283. The natural products staurosporine, quercetin, K252a and K252b have been shown to be IKK2 inhibitors, see Peet, G. W. and Li, J. J. Biol. Chem., 274, 32655-32661 (1999) and Wisniewski, D., et al., Analytical Biochem. 274, 220-228 (1999). Synthetic inhibitors of IKK2 have also been described, see Burke, et al. J. Biol. Chem., 278, 1450-1456 (2003), Murata, et al., Bioorg. Med. Chem. Lett., 13, 913-198 (2003), Murata, et al., Bioorg. Med. Chem. Lett., 14, 4013-4017 (2004), and Murata, et al., Bioorg. Med. Chem. Lett., 14, 4019-4022 (2004) have described IKK2 inhibitors.\nThus, attempts have been made to prepare compounds that inhibit IKK2 activity and a number of such compounds have been disclosed in the art. However, in view of the number of pathological responses that are mediated by IKK2, there remains a continuing need for inhibitors of IKK2 which can be used in the treatment of a variety of conditions.\nThe present inventors have discovered novel indole carboxamide compounds, which are inhibitors of kinase activity, in particular inappropriate IKK2 activity. Such indole carboxamide derivatives are therefore useful in the treatment of disorders associated with inappropriate kinase, in particular inappropriate IKK2 activity in particular in the treatment and prevention of disease states mediated by IKK2 mechanisms including inflammatory and tissue repair disorders, particularly rheumatoid arthritis, inflammatory bowel disease, asthma, rhinitis, and COPD (chronic obstructive pulmonary disease); osteoarthritis, osteoporosis and fibrotic diseases; dermatosis, including psoriasis, atopic dermatitis and ultraviolet radiation (UV)-induced skin damage; autoimmune diseases including systemic lupus erythematosus, multiple sclerosis, psoriatic arthritis, alkylosing spondylitis, tissue and organ rejection, Alzheimer's disease, stroke, atherosclerosis, restonosis, diabetes, glomerulonephritis, cancer, including Hodgkins disease, cachexia, inflammation associated with infection and certain viral infections, including acquired immune deficiency syndrome (AIDS), adult respiratory distress syndrome, and Ataxia Telangiestasia."}
{"text": "There are numerous diseases and disorders that can effect corneal tissue and which can, as a result, adversely effect or eliminate vision. For example, allergies, conjunctivitis, corneal infections, Fuchs' dystrophy (deterioration of corneal endothelial cells), varicella-zoster virus, iridocomeal endothelial syndrome, keratoconus, ocular herpes and a number of other conditions, as well as congenital corneal abnormalities can be responsible for corneal damage or irregularity that may affect vision. In an endeavour to restore sight or improve vision in people suffering from corneal abnormalities it has become particularly common to perform corneal transplant operations where the abnormal corneal tissue is removed and replaced using fine sutures with normal corneal tissue obtained from a donor. Although corneal transplant operations enjoy a high rate of success there are nonetheless some problems which can occur, such as rejection of the replacement cornea and ocular fibrosis or scarring. Even in the case of a successful corneal transplantation it is necessary for subsequent administration of immunomodulatory agents. Non-compliance by the patient with prescribed dosing regime of immunomodulating agents may give rise to tissue rejection. There is, accordingly, a need for improved means of prolonging corneal graft survival and preventing tissue rejection as well as for the provision of approaches for therapy of ocular infection, wounds and fibrosis and for therapy of other ocular disorders, for example.\nThe cornea is a highly organised group of cells and proteins which unlike most tissue is clear, and does not contain blood vessels to nourish or protect against infection. The cornea receives nutritional supply from tears and the aqueous humor found in the anterior chamber located behind it. The cornea is composed of five basic layers, namely the protective external epithelium, Bowman's layer which is located below the epithelial basement membrane and is composed of collagen fibers, the stroma which consists primarily of water and collagen and is located beneath Bowman's layer; and Descemet's membrane located beneath the stroma, which is composed of collagen fibers produced by the endothelial cells located in the lower endothelium. The endothelial cells are essential in maintaining clarity of the cornea by removing excess fluid from the stroma.\nIrreversible immunological rejection is the major cause of human corneal graft failure (1), despite the immunologically privileged nature of the eye (2). The histological correlates of rejection include local upregulation of major histocompatability complex and adhesion molecules, an influx of mononuclear cells into the cornea and anterior chamber, and local production of some inflammatory cytokines (3-7). The major target of corneal graft rejection is the corneal endothelium. Human (but not rodent) corneal endothelium is essentially amitotic (8), so that damage to the monolayer during graft rejection cannot be repaired.\nGene therapy has the potential to influence an allograft response through local expression of a modulatory gene product within transplanted donor tissue. The present inventors consider that the cornea may be uniquely amenable to such an approach because of its small size, which may allow modification of the whole tissue, and because of the ease with which a donor cornea may be manipulated in vitro and stored for a considerable period (for example, up to 28 days) prior to transplantation. The anatomical location and clarity of the cornea allow in vivo assessment of the entire graft in the post-operative period and loss of function is easy to detect. Furthermore, the cornea and anterior chamber are at least partially immunologically privileged sites (2), which may allow the use of otherwise immunogenic or pro-inflammatory vectors."}
{"text": "The ability to provide for the fuel and energy needs of the world's growing population has emerged as one of the great challenges of this century. Current fuel and energy needs are primarily met by non-renewable fossil fuels, a source that is both unsustainable and increasingly cost-inefficient. Therefore, new approaches to solving the world's energy needs are required to address these mounting concerns.\nAmong forms of plant biomass, lignocellulosic biomass is particularly well-suited for energy applications because of its large-scale availability, low cost, and environmentally benign production. In particular, many energy production and utilization cycles based on cellulosic biomass have very low greenhouse gas emissions on a life-cycle basis. The primary obstacle impeding the more widespread production of energy from biomass feedstocks is the general absence of low-cost technology for overcoming the recalcitrance of biomass feedstocks to conversion into useful products. Lignocellulosic biomass contains carbohydrate fractions (e.g., cellulose and hemicellulose) that can be converted into ethanol or other end-products including lactic acid and acetic acid. In order to convert the carbohydrate fractions, the cellulose or hemicellulose must ultimately be converted, or hydrolyzed, into monosaccharides. The hydrolysis of cellulose and hemicellulose has historically proven to be problematic.\nCellulose digesting anaerobic bacteria are of great potential utility because they can be used to produce ethanol or other fuels from abundant substrates such as forestry, municipal and agricultural waste. However, it has been challenging to realize the potential utility of biomass because of difficulty in the genetic manipulation of anaerobic bacteria and a lack of understanding of their metabolic biochemistry. Genome sequence data and recent advances in biotechnological tools for genetic modification of Clostridium thermocellum and other similar organisms have made it possible to make progress in the utilization of biomass for fuel, but the great complexity of metabolism makes it difficult to achieve a desired outcome such as near theoretical ethanol yield from cellulosic substrates.\nMany microorganisms can metabolize glucose, cellulose or cellodextrins anaerobically, but they vary in the pathways utilized and the products generated. It has been demonstrated in genetically modified Thermoanaerobacterium saccharolyticum that glucose and cellobiose can be fermented to ethanol at very close to theoretical yield, but similar genetic manipulations in Clostridium thermocellum have not had the same outcome. Argyros et al. “High ethanol titers from cellulose using metabolically engineered thermophilic, anaerobic microbes.” Appl. Env. Microbiol., 2011 doi:10.1128/AEM.00646-11 (epub ahead of publication).\nClostridium thermocellum has both cellulolytic and ethanologenic fermentation capabilities and can directly convert a cellulose-based substrate into ethanol. However, C. thermocellum possesses a branched carbon utilization pathway that generates undesirable products, and thus its yield of ethanol is low. Furthermore, C. thermocellum is not as amenable to manipulation for ethanol production as T. saccharolyticum. The difficulty in manipulating C. thermocellum for ethanol production is exemplified more clearly when the carbon utilization pathways from C. thermocellum and T. saccharolyticum are compared. In homoethanologenic T. saccharolyticum, the carbon atoms from glucose flow down a linear central metabolic pathway to ethanol (FIG. 1A). In C. thermocellum, a different set of enzymes is present and thus the carbon utilization pathway (FIG. 1B) is different that the carbon utilization pathway in T. saccharolyticum. The difference in the carbon-utilization pathways of C. thermocellum compared to T. saccharolyticum makes it infeasible to produce ethanol at theoretical yield with the same modifications.\nMany enzymes in carbon-utilizing metabolic processes use a nicotinamide adenine dinucleotide as a cofactor. There are two common types of nicotinamide adenine dinucleotide cofactors, NAD+ and NADP+. Each can exist in a reduced or oxidized form. In order to maintain steady state, each cofactor involved in a reaction must be regenerated at the same rate it is consumed. In other words, the cell must be reduction-oxidation (“redox”) balanced. Enzymes are typically specific for (i.e. react with) either the phosphorylated (NADP+, NADPH) or non-phosphorylated (NAD+, NADH) nicotinamide cofactors. The specificity of an enzyme can sometimes be switched from one nicotinamide cofactor to the other by mutations in the cofactor binding region of the protein. It is also possible to find different isoforms of an enzyme that carry out the same enzymatic activity, but use different cofactors (e.g. NAD+ instead of NADP+). Isoforms with altered cofactor specificity may be found for example in different species.\nThe T. saccharolyticum oxidation-reduction reactions in the metabolic pathway from cellobiose to ethanol are:\n(1) D-glyceraldehyde 3-phosphate+phosphate+NAD+=3-phospho-D-glyceroyl phosphate+NADH+H+ (catalyzed by glyceraldehyde-3-phosphate dehydrogenase)\n(2) pyruvate+CoA+oxidized ferredoxin=acetyl-CoA+CO2+reduced ferredoxin+H+ (catalyzed by pyruvate oxidoreductase)\n(3) reduced ferredoxin+NADH+2 NADP++H+=oxidized ferredoxin+NAD++2 NADPH (catalyzed by NADH-dependent reduced ferredoxin:NADP+oxidoreductase)\n(4) acetyl-CoA+NADPH+H+=acetaldehyde+CoA+NADP+ (catalyzed by acetaldehyde dehydrogenase)\n(5) acetaldehyde+NADPH+H+=ethanol+NADP+ (catalyzed by alcohol dehydrogenase)\nReactions 1-5 above are redox and cofactor balanced. A single polypeptide called AdhE contains both catalytic activities of steps 4 and 5. Activity of AdhE is detectable with both NADH and NADPH cofactors (See Shaw et al., “Metabolic engineering of a thermophilic bacterium to produce ethanol at high yield.” PNAS 2008. 105(37): 13769-74). In C. thermocellum, activity can be detected for both cofactors in the alcohol dehydrogenase reaction, but the aldehyde dehydrogenase reaction is specific to NADH only (See Brown et al., “Mutant alcohol dehydrogenase leads to improved ethanol tolerance in Clostridium thermocellum.” PNAS 2011. 108(33): 13753-7 and Rydzak et al., “Growth phase-dependent enzyme profile of pyruvate catabolism and end-product formation in Clostridium thermocellum ATCC 27405.” J. of Biotech. 2009. 104(3-4): 169-75). Therefore, reaction 4 above cannot occur in C. thermocellum. Reaction 4 can occur with NADH as the cofactor, but use of NADH would lead to an overabundance of NADPH and depletion of NADH in the cell. The oxidation-reduction reactions in C. thermocellum in the pathway from cellobiose to ethanol are the same as 1-5 above, but with the addition of two more:\n(6) oxaloacetate+NADH+H+=malate+NAD+ (catalyzed by malate dehydrogenase)\n(7) malate+NADP+=pyruvate+CO2+NADPH (catalyzed by malic enzyme)\nThe net effect of these two additional reactions in C. thermocellum is that electrons are transferred from NADH to NADPH. This leads to a further accumulation of NADPH and makes the pathway from cellobiose to ethanol unbalanced for cofactors and therefore infeasible in this configuration. As a result, C. thermocellum strains lacking the ability to make other end products (e.g. mutants for lactate dehydrogenase and phosphotransacetylase) show poor ethanol productivity and secrete amino acids that consume NADPH during their biosynthesis.\nConsequently, in order to optimize ethanol production in C. thermocellum, there is a need for mutant strains of C. thermocellum that are reduction-oxidation and cofactor balanced.\nThe present invention relates to cellulose-digesting organisms that have been genetically modified to allow the production of ethanol at high yield by changing cofactor usage and/or production at key steps of central metabolism."}
{"text": "1. Field of the Invention\nThis invention is related to a system, apparatus and method for the non-lethal restraint of a vehicle through the use of an entanglement device that will entangle such vehicle. The non-lethal entanglement device incorporates a plurality of tendrils, filaments, tentacles or straps, or a combination there of, that are propelled from a housing by compressed gas or by pressure generated by a gas generator of the type commonly used in air bag deployment apparatus. Filaments that are launched from the device may be attached to frangible balls carrying adhesive substances or barbed capture elements that will adhere or stick to a target surface. The filaments are designed to assist in entangling a target vehicle and restrain the entangled target.\n2. Description of Related Art\nDevices for stopping a fleeing vehicle include barriers, tire spike strips, caltrops, snares and electrical system disabling devices. The well-known spike strip, with spikes projecting upwardly from a base structure, usually an elongated structure, either a rolled up device or an accordion type device, is designed to be placed on a road in anticipation that a fleeing target vehicle will be driven over the spike strip. Once tire contact is made with the spike strip it will cause deflation of vehicle tires, and eventually cause the vehicle to stop due to the vehicle being difficult to control on flat tires. Barriers, such as concrete barricades, can be effective but to set them up is time consuming and barriers are not particularly portable. Barriers are normally used only in static, non-portable placements. Caltrops are small pyramid-like devices with four projections at about one hundred nine degrees offset to each other. When a caltrop is deployed it will have one of the projections, typically a spike, pointing vertically upward. A vehicle will have to drive directly over a caltrop so as to puncture a tire. Snares, such as a net that is stretched across a road, can also be effective in stopping a vehicle but such snares are time consuming to set up, are generally bulky and heavy and need to be anchored securely to restrain a moving vehicle.\nIt is also known that there are vehicle immobilization devices that depend on an electrical charge to disable the vehicle. These devices may work with vehicles that have electronic ignition, on-board engine and component control computers, microchip controlled systems, and other electrical control components that can be destroyed or made inoperative through the application of a voltage surge or other high energy electrical pulse or charge provided to a vulnerable electrical systems. Such electrical based-disabling devices may have limited or no effect on vehicles with shielded electrical systems or on vehicles that are not electronically dependent to operate.\nRegardless of the type of vehicle stopping device used it is enlightening to realize that the energy needed to stop a fast moving vehicle is significant. Kinetic energy attained by a moving vehicle is on the order of one-half of the mass of the vehicle times the velocity squared. Thus a large truck, such as the well known Humvee (High Mobility Multipurpose Wheeled Vehicle (“HMMWM”), can exhibit more than ten times the kinetic energy of a small passenger car. Thus it is important to realize that a large amount of energy is needed to stop a moving vehicle and that any immobilization device will need to be effective against an array of vehicles to be capable of stopping heavy as well as light vehicles. In this invention the undercarriage immobilization device described herein is not intended to go “head-to-head” with the fleeing vehicle, but is designed to use motion of the vehicle components to assist in the restraint or capture of a target vehicle and thus render the vehicle immobile or significantly slowed to the point of being easily overtaken by law enforcement personnel.\nTo reduce the complexity and length of the Detailed Specification, and to fully establish the state of the art in certain areas of technology, Applicants herein expressly incorporate by reference material identified in the following publications:    National Institute of Justice, “Department of Defense Non-lethal Weapons and Equipment Review: A Research Guide for Civil Law Enforcement and Corrections,” October 2004, NCJ-205293.    Greg Lucas, “Bay Area's New Efforts in the War on Terror Coast Guard Weapon: High-tech net to keep boats from off-limits areas,” San Francisco Chronicle Article, Aug. 10, 2005. (Available on the Internet.).    MIL-HDBK 5 “Military Handbook Metallic Materials and Elements for Aerospace Vehicle Structures,” United States DOD, Dec. 1, 1998.    MIL-HDBK 17 “Military Handbook Composite Materials Handbook,” United States DOD, Dec. 12, 2002.    Honeywell Spectra Technical Bulletin, HON-PF-PS10, (Available on the Internet).    Steven H. Scott, “Sticky Foam as a Less-Than-Lethal Technology,” Sandia National Laboratory, US DOE Contract No DE-AC04-96AL8500, CIRCA 1994.    T. D. Goolby and K. J. Padilla, “Sticky Foam Restraining Effectiveness Human Subject Tests for the Less-Than-Lethal Foam Project,” Sandia Report, Jul. 8, 1994 UNCI (Available on the Internet).\nThe applicants believe that the material incorporated above is “non-essential” in accordance with 37 CFR 1.57, because it is referred to for purposes of indicating the background of the invention or illustrating the state of the art. However, if the Examiner believes that any of the above-incorporated material constitutes “essential material” within the meaning of 37 CFR 1.57(c)(1)-(3), applicants will amend the specification to expressly recite the essential material that is incorporated by reference as allowed by the applicable rules."}
{"text": "The present invention relates to rotary sprinklers, and is particularly directed to an improvement to the type of rotary sprinkler described in U.S. Pat. No. 4,261,515.\nThe above Application describes a rotary sprinkler comprising a nozzle having an inlet connectable to a source of pressurized water and formed with an axial bore through which the water exits in the form of an axial jet, a rotor in the path of the axial jet, and means for floatingly mounting the rotor for axial, lateral and rotary movements with respect to the nozzle bore, the underface of the rotor being formed with at least one groove (preferably two) extending from the rotor center to its outer edge for deflecting the jet laterally of the sprinkler and for imparting a rotory motion to the rotor."}
{"text": "The present invention relates to a dental handpiece for the polymerization of photosetting compounds or resins compatible with the power supply of other handpieces.\nThe invention refers in particular to an instrument employed by dental surgeons for the polymerization of photosetting compounds or resins utilised especially in dental prostheses,\nPhotosetting resins are applied in a semifluid state to patients\"\" dental prostheses and are then hardened by a polymerization process. Such a polymerization process involves heating the resins for a certain length of time, using a light source that emits blue light, i.e. having an emission spectrum with a wavelength centered around 470 nm.\nHandpieces according to the prior art currently use halogen lamps as their light source. The light coming from the halogen lamps is filtered by dichroic filters, so that a blue light having a light spectrum with a wavelength between 430 and 510 nm is obtained. The light emerging from the dichroic filters is transported, through optic fibres, to the outlet of the dental handpiece, enabling the dentist to direct it onto the resin to be polymerized, applied to the patient\"\"s dental prosthesis.\nThe known dental handpieces using halogen lamps as polymerizing lamps have various drawbacks.\nIn order to achieve polymerization of the resin, the light emerging from the dental handpiece must have a power of approximately 500 mW. Since halogen lamps have a very low power efficiency (efficiency value xcex7=0.5-1%), this results in a high energy dissipation by the polymizering lamps, which reaches values between 50 and 100 W.\nFor this reason, handpieces for the polymerization of photosetting compounds must have a special power supply, able to supply the polymerizing lamps with a power that ranges from 50 to 100 W. Consequently, the seats in which the patients sit, normally known as dental chair units, provide a connector element connected to the power supply and specifically designed for the polymerizing handpiece.\nFurthermore, with the output power that can be obtained from halogen lamps, rather long resin polymerization times are required; this can be very tiresome both for the patient undergoing the procedure and for the dental surgeon who has to perform it.\nIn an attempt to reduce resin polymerization times, handpieces with different light sources have been put on the market. Handpieces are per se known that use as their light source gas-discharge lamps (plasma torches) that emit high-power white light, which is filtered in order to obtain a light beam with an emission spectrum centred on blue. Also known are dental handpieces using laser as their light source, which directly transmit a light beam with an emission spectrum centred on blue.\nWith this type of dental handpiece it is possible to obtain a light beam with a power output approximately ten times greater compared to that obtained with halogen lamps, therefore shorter polymerization times are achieved. However, such dental handpieces, apart from being very expensive, have the problem of needing to dissipate a large amount of heat and inevitably they must provide a heat dissipation or cooling system, making them excessively bulky and heavy. Furthermore, such heat dissipation systems often include fans driven by motors which must be powered by the dedicated electrical supply of the polymerizing handpiece.\nAlso known, especially at an experimental level, are handpieces for polymerization of photosetting compounds using as their light source LEDs that emit in the blue range.\nAll known handpieces using halogen lamps, plasma torches, laser or LEDs as their light source, have an independent electrical supply system, dedicated soleley to the electrical supply of said light sources. This makes the dental handpiece bulky and lacking in versatility because of the purpose-built circuitry provided for such light sources and also the need to provide the dental chair unit with special connectors for the power supply of such handpieces.\nThe object of the invention is to eliminate said drawbacks, providing a dental handpiece for polymerization of photosetting compounds or resins that is compatible with the power supply of other handpieces.\nA further object of the present invention is to provide a handpiece of reduced dimensions and bulk that is versatile, economical and easy to produce.\nThe handpiece according to the invention comprises a connector element shaped so as to be able to connect to external electrical connectors utilised for other types of handpieces, such as for example the handpiece or ultrasound scaler for removal of plaque and caries, or the dental turbine for cleaning and washing dental prostheses.\nThis is made possible thanks to the fact that in the handpiece according to the invention LEDs that emit in the blue range are used as the light source. In fact LEDs require a decidedly lower power supply than halogen lamps, plasma torches or laser. This has made it possible to take the electrical supply for bias of the LEDs directly from the electrical supply circuit provided in the external connector of other types of handpieces.\nIn the case of connection of the handpiece for polymerization according to the invention to the connector of an ultrasound scaler, the electrical energy for bias of the LEDs is taken from the electrical circuit provided in the connector of the scaler and destined to supply the ultrasound transducer of the scaler.\nFor this purpose, in the handpiece according to the invention, an equivalent resonator circuit is provided that must tune the resonance frequency of the supply circuit provided in the external connector of the scaler. Downstream of the equivalent resonance circuit and upstream of the LED assembly of the polymerizing handpiece, a voltage and current stabilizer or regulator is provided in order to generate a constant voltage and direct current signal for correct bias of the LEDs.\nDental turbines include a light source to illuminate the work surface. For the electrical supply to this light source of the dental turbine, the external connectors of dental turbines carry an electrical supply signal taken from an electrical supply source. In the case of connection of the handpiece according to the invention to an external connector for dental turbines, the supply of the LEDs for polymerization is taken from the electrical supply carried by the external connector of the turbine.\nIn this case the voltage supplied by the external connector is too low to be able to supply the LEDs, therefore an electronic voltage booster or step-up transformer circuit capable of providing an adequate voltage to bias the LEDs is provided inside the handpiece according to the invention. Downstream of the step-up transformer circuit a voltage and current stabilizing circuit is provided to maintain a constant voltage and direct current during power supply of the LEDs.\nThe advantages of the handpiece for the polymerization of photosetting compounds according to the invention appear evident. Since it can adapt to the connectors used for the power supply of handpieces of other types, it allows the connectors and circuitry specifically dedicated to its supply to be eliminated from the dental chair unit, with considerable economic advantages as regards size and complexity."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of electrical connectors. More particularly, this invention relates to an electrical connector commonly used for audio and video applications which is commonly referred to as an audio jack connector, RCA connector or phono connector.\n2. Background of the Invention\nThe term \"audio jack connectors\", \"phono connectors\", \"RCA connectors\" and the like are used synonymously herein. These terms are intended to generically embrace coaxial connectors which are commonly referred to by these names in the electronics industry as well as other coaxial connector devices which operate in a similar manner. These terms are not intended to restrict the usage of the connector to audio or phonograph applications since they are commonly used for other purposes including video, radio frequency (RF) and intermediate frequency (IF) applications.\nConnectors such as those described above are commonly used for many electrical and electronic applications. These connectors are particularly popular in consumer electronics devices. As such, large quantities of these connectors are in use. In addition, there is a great deal of demand for such devices to be produced and installed in electronic systems at low cost to the manufacturer. Similar coaxial connection devices are also commonly used in radio frequency applications and the like which can benefit from decreased cost.\nTraditionally, such connectors are installed manually by either soldering a pair of solder tabs to wires or cables which are connected to other electrical circuitry. Mechanical support is obtained by either a solder connection or by passing the connector through a hole in an electrical enclosure and securing the connector with a threaded nut and washer or similar connection procedures.\nObviously, such manufacturing techniques are labor intensive and result in higher product cost and lower reliability. It is thus desirable to provide a system for producing such connectors in a manner which is more conducive to automation in the manufacture of the end user product. In addition, existing connectors are made of several metal parts and insulators which are assembled together. It is also desirable to reduce the number of parts in the assembly to simplify the design and reduce or eliminate the possibility of assembly errors.\nThese and other drawbacks of the prior art connectors are ameliorated by the improved coaxial connector of the present invention."}
{"text": "In general, conventional rate converters include two major computation blocks, as illustrated in FIG. 1. A first block is a conversion rate tracking loop that determines the ratio between an input sample rate and output sample rate. Once the ratio is determined, the conversion rate tracking loop generates the corresponding input sample index for each output sample.\nThe second computation block is the sample interpolator. The function of this block is to interpolate the input data sample and to generate the output data sample with the real value index. One problem with conventional rate converters is that they create aliasing from the resampling, and it is relatively difficult to produce an output signal having a high Signal-to-Noise Ratio (SNR) and having a relatively flat frequency response within standard frequency ranges for audio signals.\nEmbodiments of the invention address this and other limitations of the prior art."}
{"text": "This invention relates generally to a fluid valve and more particularly to a switching logic element having a normally closed characteristic.\nSwitching logic elements utilize an unbalanced poppet which is spring loaded on a seat; however, without additional fluid pressure to bias the poppet against the seat, flow can occur between the input and output ports. These valves are called \"normally open\" because of the need for fluid pressure to close them. The complementary logic element--normally closed--has not previously been available. Such a normally closed logic element would require a positive or pilot pressure input to be opened and to remain open.\nThe need for a pilot pressure is an operating requirement for fail safe applications, i.e. if the source of fluid pressure is cut off either inadvertently or because of system shut down, the valve should close to prevent loss of load supporting pressure by flow through the logic element."}
{"text": "1. Field of the Invention\nThe present invention relates to a powered vehicle door closing system and specifically to a powered vehicle door closing system suitable for an automotive vehicle such as a van with a sliding door moveable between open and closed positions relative to a vehicle body opening, and more specifically to a system which is capable of forcibly and automatically moving a latch member employed in a lock unit from a half-latched position (a nearly-closed position of the sliding door) to a fully-latched position (a fully-closed position of the sliding door) by powering the final, low-displacement/high-force movement of the sliding door.\n2. Description of the Prior Art\nRecently, there have been proposed and developed various powered vehicle door closing systems which can automatically move a latch member from a half-latched position to a fully-latched position. One such powered vehicle door closing system has been disclosed in Japanese Patent Provisional Publication (Tokkai Heisei) No. 1-105886. The powered door closing system disclosed in the Japanese Patent Provisional Publication No. 1-105886 is applied to a door lock for an automobile sliding door. The prior art door closing system has three switches, namely a first switch for detection of a half-latched state of the latch member, a second switch for detection of a fully-latched state of the latch member, and a third member for detection of a stand-by position of a moveable drive lever (a portion of a force-transmitting linkage) by way of which the latch member can be shifted from the half-latched position to the fully-latched position. The first switch consists of a pair of electrical contacts, one being a stationary electrical contact provided in the vehicle body and the other being a spring-loaded, plunger-type electrical contact provided onto the door for contact with the stationary contact upon shift to the half-latched position of the latch member via the manual door operation. The first switch is responsive to the movement of the sliding door in such a manner as to rotate the drive lever away from its stand-by position by way of normal rotation (positive rotation) of a drive motor such as a reversible electric motor when the sliding door reaches the half-latched position of the latch member, and as a result the latch member is forcibly moved to its fully-latched position. The second switch is responsive to the movement of the latch member in such a manner as to rotate the drive lever toward the stand-by position by way of reverse-rotation (negative rotation) of the drive motor when the latch member reaches the fully-latched position. The third switch is responsive to the movement of the drive lever in such a manner as to stop the drive motor and consequently to maintain the drive lever at the stand-by position immediately when the drive lever reaches the stand-by position. Each of the second and third switches consists of an ordinary limit switch which is so designed to cut off power automatically at or near the limit of travel of a moving object controlled by electrical means. The conventional system requires an accurate installation of the first, second, and third switches, since respective spring-loaded switching contacts or points of the second and third switches are mechanically operated and additionally the stationary electrical contact (included in the first switch) is provided in the vehicle body, while the other plunger-type electrical contact is provided onto the door for contact with the stationary contact. When installing the three switches on the vehicle, the conventional system requires a complicated wiring harness. This results in an increased assembling time of the respective switches onto the vehicle. Due to a complicated, costly structure of the conventional system, total manufacturing costs of the automotive vehicle with an auto door closing system is increased. Owing to an inherent switching characteristic of the limit switch with the spring-loaded mechanical contact, there is a possibility that a switched-ON operation of the limit switch cannot be completed, particularly during manual quick door closing operation with great momentum. There is a possibility of malfunction of the system."}
{"text": "Refrigerator appliances generally include a cabinet that defines a chilled chamber for receipt of food items for storage. For example, the cabinet can define a fresh food chamber and a freezer chamber. The fresh food chamber can be maintained at a temperature greater than the freezing point of water. Conversely, the freezer chamber can be maintained at a temperature equal to or less than the freezing point of water.\nRefrigerator appliances generally also include one or more shelf assemblies positioned within the fresh food chamber and/or the freezer chamber to facilitate storage and/or organization of any food items positioned therein. Typically, the shelf assemblies are adjustable in height using, e.g., a cantilevered track assembly. Accordingly, with such a configuration, the user may customize the fresh food chamber and/or the freezer chamber of the refrigerator appliance to meet their specific needs.\nModern refrigerator appliances can also include lighting systems attached to or embedded within the one or more shelf assemblies. The lighting systems can, e.g., illuminate the shelf assembly itself, or alternatively can illuminate a lower shelf assembly. However, in order for such refrigerator appliances to provide such lighting systems with electrical power, one or more electric wires are generally required to be disconnected and reconnected as the shelf assembly is adjusted. Certain consumers may find such electrical wires unsightly and may find the additional steps of disconnecting and reconnecting wires undesirable.\nAccordingly, a refrigerator appliance having one or more shelf assemblies capable of connecting a lighting system to an electrical source without requiring disconnection and reconnection of electrical wires would be useful. Moreover, a refrigerator appliance having one or more shelf assemblies capable of connecting a lighting system to an electrical source without any visible connections would be particularly beneficial."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical information recording medium being utilized for a recording and reproducing apparatus or a drive, which records an information in the optical information recording medium with moving the medium relatively and reads out the information from the medium, and a substrate and a manufacturing method for the optical information recording medium, particularly, relates to an optical information recording medium, which can be recorded and reproduced in a high density and a large capacity, and a substrate and a manufacturing method for the optical recording medium.\n2. Description of the Related Art\nCurrently, there existed a system, which can optically record and reproduce information by using a disk shaped medium, as a system of an information recording medium for reading out information with moving the medium relatively. A disk system medium can be divided into several types such as a read only type or a ROM (read only memory) type, a recordable or a write once type or an “R” type and a rewritable type or a RW type. Generally, a recording density of a medium is high in a ROM type and low in “R” and RW types. For instance, in a DVD (digital versatile disk) system, which utilizes a laser beam having a wavelength of 635 through 650 nm, introduced into the market in 1996, a read only type disk such as a DVD-ROM disk and a DVD-Video disk was introduced first. A recording capacity of a read only type DVD disk is 4.7 GB. On the other hand, a recording capacity of a rewritable type DVD disk such as a DVD-RAM (DVD-Random Access Memory) disk is 2.6 GB, that is, the recording capacity of a DVD-RAM disk is almost 55% of that of a DVD-ROM disk. Researches and developments for increasing a recording capacity of a rewritable type disk are progressing. However, a system having a same capacity as that of a DVD-ROM disk has not been developed yet.\nIn the case of a rewritable type disk, a recording format on a disk and materials of recording medium are key technologies. In a DVD-RAM disk, a land-groove recording method, which is utilized for recording information on both land and groove of an optical information recording medium, has been used. In this method, an address necessary to recording and reproducing is recorded with cutting a land and a groove at each specific period of time.\nFIG. 19 is a fragmentary plan view of a disk partially enlarging a microscopic construction 20, that is, a physical format construction of the land-groove recording method disk. In FIG. 19, the drawing shows an external appearance of the construction when the disk is not recorded and grooves 21a through 21c, hereinafter represented by the groove 21, are formed in parallel to each other. Lands 22a and 22b, hereinafter represented by the land 22, are allocated in between the grooves. An information is recorded on both the groove 21 and the land 22 when recording. A plurality of address pits 23a through 23n, which is necessary to recording and reproducing, is formed by cutting the groove 21 and land 22. The addresses occupy an area 24 in conjunction with an accompanied signal, so that the area 24 prevents a total capacity of a disk from increasing. In other words, a limited area of a disk can not be effectively utilized because of the area 24.\nFurther, with respect to recording material, in consideration of interchangeability with a read only type DVD disk drive, a phase change recording method, which does not utilize a magnetic head, is suitable. However, this method is defective in reflectivity, which is much inferior to a read only type disk or a write once type disk of utilizing dye. Accordingly, the low reflectivity results in that a recording capacity can not be increased.\nBy combining a microscopic construction or a physical format construction with a phase change material for high density recording and by optimizing them, a recording capacity equivalent to that of a read only type DVD disk can be realized.\nAn optical disk, which records addresses with scattering over a disk without having an inherent address area such as the area 24 shown in FIG. 19, is considered as a format for a large capacity optical disk. In other words, an address area such as the area 24 is not provided, so that a recording capacity can be increased as many as that of a DVD-ROM disk. However, in a case that a main recording signal is recorded in a vicinity area of an address signal, the main recording signal may interfere in the address signal and an error may occur by this method, and then rewriting can not be performed any more. Conversely, the address signal may leak into and interfere in the main recording signal and result in a reading out error."}
{"text": "This invention relates generally to an apparatus and method for an improved nanocalorimeter, and more specifically, to a system and method for an improved nanocalorimeter for measuring the heat released or absorbed during chemical reactions.\nCalorimetry is used to measure enthalpic changes, including enthaplic changes arising from reactions, phase changes, changes in molecular conformation, temperature variations, and other variations of interest that may occur for a particular specimen. By measuring enthalpic changes over a series of conditions, other thermodynamic variables may be deduced. For example, measurements of enthalpy as a function of temperature reveal the heat capacity of a specimen, and titrations of reacting components can be used to deduce the binding constant and effective stoichiometry for a reaction. Calorimetry measurements are useful in a broad variety of applications, including, for example, pharmaceuticals (drug discovery, decomposition reactions, crystallization measurements), biology (cell metabolism, drug interactions, fermentation, photosynthesis), catalysts (biological, organic, or inorganic), electrochemical reactions (such as in batteries or fuel cells), and polymer synthesis and characterization, to name a few. In general, calorimetry measurements can be useful in the discovery and development of new chemicals and materials of many types, as well as in the monitoring of chemical processes. Standard calorimeters require relatively large samples (typically about 0.5 ml to 10 liters) and usually measure one sample at a time. As such, these systems cannot be used to measure very small samples, as might be desired for precious or highly reactive materials. Furthermore, standard calorimeters cannot be used effectively to monitor a large number of reactions of small sample size in parallel, as is required in order to perform studies using combinatorial chemistry techniques.\nIn recent years, researchers and companies have turned to combinatorial methods and techniques for discovering and developing new compounds, materials, and chemistries. For example, pharmaceutical researchers have turned to combinatorial libraries as sources of new lead compounds for drug discovery. As another example, Symyx Technologies® is applying combinatorial techniques to materials discovery in the life sciences, chemical, and electronics industries. Consequently, there is a need for tools that can measure reactions and interactions of large numbers of small samples in parallel, consistent with the needs of combinatorial discovery techniques. Preferably, users desire that these tools enable inexpensive measurements and minimize contamination and cross-contamination problems.\nIn some cases, the sample to be studied is precious, and it might not be acceptable to use the relatively large amount of material required by a standard microcalorimeter to perform only one measurement. For example, one may desire to study a natural extract or synthesized compound for biological interactions, but in some cases the available amount of material at concentrations large enough for calorimetry might be no more than a few milliliters. Performing a measurement in standard microcalorimeters, such as those sold, for example, by MicroCal® Inc. (model VP-ITC) or Calorimetry Sciences Corporation® (model CSC4500), requires about 1 ml of sample, which means that one would possibly be Faced with using a majority or all of the precious material for one or a small series of measurements. Tools that enable calorimetric measurements with much smaller sample sizes would be helpful in overcoming this limitation.\nOne of the most popular uses of combinatorial techniques to date has been in pharmaceutical research. Pharmaceutical researchers have turned to combinatorial libraries as sources of new lead compounds for drug discovery. A combinatorial library is a collection of chemical compounds which have been generated, by either chemical synthesis or biological synthesis, by combining a number of chemical “building blocks” as reagents. For example, a combinatorial polypeptide library is formed by combining a set of amino acids in every possible way for a given compound length (i.e., the number of amino acids in a polypeptide compound). Millions of chemical compounds can theoretically be synthesized through such combinatorial mixing of chemical building blocks.\nOnce a library has been constructed, it must be screened to identify compounds, which possess some kind of biological or pharmacological activity. For example, screening can be done with a specific biological compound, often referred to as a target, that participates in a known biological pathway or is involved in some regulation function. The library compounds that are found to react with the targets are candidates for affecting the biological activity of the target, and hence a candidate for a therapeutic agent.\nThrough the years, the pharmaceutical industry has increasingly relied on high throughput screening (HTS) of libraries of chemical compounds to find drug candidates. HTS describes a method where many discrete compounds are tested in parallel so that large numbers of test compounds are screened for biological activity simultaneously or nearly simultaneously. Currently, the most widely established techniques utilize 96-well microtitre plates. In this format, 96 independent tests are performed simultaneously on a single 8 cm×12 cm plastic plate that contains 96 reaction wells. These wells typically require assay volumes that range from 50 to 500 μl. In addition to the plates, many instruments, materials, pipettors, robotics, plate washers and plate readers are commercially available to fit the 96-well format to a wide range of homogeneous and heterogeneous assays. To achieve faster testing, the industry is evolving to plates that contain 384 and 1536 wells.\nA variety of measurement approaches has been used to screen combinatorial libraries for lead compounds, one of which is the inhibitor assay. In the inhibitor assay, a marker ligand, often the natural ligand in a biological pathway, is identified that will bind well with the target protein molecule. The assay requires the chemical attachment of a fluorescent molecule to this marker ligand such that the fluorescent molecule does not affect the manner in which the marker ligand reacts with the target protein. To operate an inhibitor assay, the target protein is exposed to the test ligands in microtitre wells. After a time necessary for reaction of the test ligand to the target protein, the marker ligand is applied. After a time for reaction with the marker ligand, the wells are rinsed such that non-reacted marker ligand is washed away. In wells where the target protein and the test ligand have reacted, the test ligand blocks the active site of the target protein so the marker ligand cannot react and is washed away, while in cells where the target protein and test ligand have not reacted, the marker ligand reacts with the target protein and is not washed away. By investigating the wells for the presence of fluorescence after the washing, reactions of test ligands and target proteins can be determined as having occurred in wells where no fluorescence is observable.\nHowever, the inhibitor assay requires time and expense to develop the assay. The principal components that need development are discovering a marker ligand and attaching a fluorophore to the marker in a manner that does not affect its reaction with the target protein. Attaching the fluorescent marker can often take 3 months of development or more and cost $250 k or more once the marker ligand is identified. An assay method that avoids such assay development, such as measuring the heat of the reaction with calorimetry, would eliminate this cost and lime delay in the discovery process.\nCalorimetry measurements are commonly utilized in biophysical and biochemical studies to determine energy changes as indications of biochemical reactions in a media. Prior techniques for measurements include using electrodes, thermopiles, optical techniques, and microcalorimeters for measurements within a sampled media. There is a great interest in developing ultra-miniature microcalorimeter devices that require very small volumes of sampled media for accurate detection and measuring of biochemical reactions on, or in close proximity to, the microcalorimeter and which can be applied in a manner to quickly measure large numbers of reactions such that it can be as efficient as assays such as inhibitor assays which can be used in HTS to screen perhaps 100,000 test ligands a day.\nThe following disclosures may be relevant and/or helpful in providing an understanding of some aspect of the present invention:\nIn Plotnikov et al., U.S. Pat. No. 5,967,659 (“Ultrasensitive Differential Microcalorimeter with User-selected Gain Setting”), a differential calorimeter is disclosed that includes sample and reference cells, a thermal shield surrounding the cells, heating devices thermally coupled to the thermal shield and the cells, a temperature monitoring system, and a control system. The temperature monitoring system monitors the temperature of the shield, cell temperatures, and temperature differentials between the cells and the shield. The control system generates output signals for control of the heating devices, with a gain setting and scan rate selected by means of a user interface. Output control signals are functions of input temperature signals and the user-selected gain setting, as well as functions of input temperature signals and the user-selected scan rate using a mapping function stored in memory.\nIn Cavicchi et al., U.S. Pat. No. 6,079,873 (“Micron-scale Differential Scanning Calorimeter on a Chip”), a differential scanning microcalorimeter produced on a silicon chip enables microscopic scanning calorimetry measurements of small samples and thin films. The chip, fabricated using standard CMOS processes, includes a reference zone and a sample zone. The reference and sample zones may be at opposite ends of a suspended platform or may reside on separate platforms. Each zone is heated with an integrated polysilicon heater. A thermopile consisting of a succession of thermocouple junctions generates a voltage representing the temperature difference between the reference and sample zones.\nIn Thundat et al., U.S. Pat. No. 6,096,559 (“Micromechanical Calorimetric Sensor”), a calorimeter sensor apparatus utilizes microcantilevered spring elements for detecting thermal changes within a sample containing biomolecules which undergo chemical and biochemical reactions. The spring element includes a bimaterial layer of chemicals on a coated region on at least one surface of the microcantilever. The chemicals generate a differential thermal stress across the surface upon reaction of the chemicals with an analyte or biomolecules within the sample due to the heat of chemical reactions in the sample placed on the coated region. The thermal stress across the spring element surface creates mechanical bending of the microcantilever. The spring element has a low thermal mass to allow detection and measuring of heat transfers associated with chemical and biochemical reactions within a sample place on or near the coated region. Deflections of the cantilever are detected by a variety of detection techniques.\nIn Lieberman, U.S. Pat. No. 6,193,413 (“System and Method for an Improved Calorimeter for Determining Thermodynamic Properties of Chemical and Biological Reactions”) a microcalorimeter includes a thin amorphous membrane anchored to a frame within an environmental chamber. Thermometers and heaters are placed on one side of a thermal conduction layer mounted on the central portion of the membrane. Samples are placed on two such membranes; each sample is heated and its individual heat capacity determined. The samples are then mixed by sandwiching the two microcalorimeters together to cause a binding reaction to occur. The amount of heat liberated during the reaction is measured to determine the enthalpy of binding."}
{"text": "1. Field of the Invention\nThe present invention relates to the filed of the digital signal processing, and more specifically, to the field of phase equivocation in burst modems.\n2. Discussion of the Prior Art\nWhen the communication exists for a long period of time, it is possible to use a conventional phase-lock loop (PLL) in a receiver to recover the carrier from a continuously transmitted signal. The PLL that is specifically designed to recover the carrier is termed as a phase carrier loop.\nHowever, in a time-division multiple access (TDMA) communication system, the available spectrum is used by letting each user to have access to the whole band for a short time (traffic burst), during which time the user transmits data as fast as the user can. The user's frequency allocation is shared with the other users who have time slots allocated at other times.\nThus, in a TDMA-type burst communication systems, the signals exist for a short period of time. Consequently, there is little time in which the phase carrier loop can slew to the proper phase in order to recover the carrier from the burst signal.\nThis situation is exacerbated for a TDMA-type burst commination system if the phase of the incoming symbol is very close to the border between the plant points. If this is the case, the slicer in the phase carrier loop can not determine to which quadrant the incoming symbol point having such phase should belong because such symbol can be equally located in either of two quadrants. Such symbol point is further referred to as a symbol point having a phase equivocation, and the phase carrier loop can not lock fast enough on an incoming signal having a phase equivocation characteristic.\nWhat is needed is to design a carrier phase loop system that is capable of resolving the phase equivocation of the incoming symbols, thus placing any incoming symbol point to a certain quadrant for further processing by the loop, thus increasing the speed of locking of the carrier phase loop on the carrier."}
{"text": "A modern telecommunications system often comprises one or more switched telephone networks and one or more Internet Protocol-based packet networks. These two different types of networks are interconnected by a media gateway, which acts as a translator between the two types of networks. The media gateway enables multimedia communications, such as voice and video, over multiple transport protocols end to end.\nBecause the media gateway connects different types of networks, one of its main functions is to convert between the different transmission and coding techniques used across the different networks. For example, a Voice-over-Internet-Protocol-capable (VoIP-capable) media gateway performs the conversion between time division multiplexed voice media that originate at a switched telephone network telecommunications terminal and VoIP datagram media that is intended for an Internet Protocol network terminal, as part of a telephone conversation between two parties; of course, the media gateway has to perform the conversion in the other direction as well. Other functions that the media gateway provides are echo cancellation, tone detection, tone generation (e.g., dual tone multi-frequency tones, etc.), and conferencing.\nSince a packet stream that is received from the Internet Protocol network comprises data packets and control packets, which contain addressing information, the VoIP media gateway converts the received packets to a time division multiplexed stream while processing the control packets. The media gateway must perform the conversion in a timely manner to minimize the possibility of packet loss, which the listening party on a call might perceive. To handle all of the packets responsively and without unacceptable delay or jitter, the media gateway uses digital signal processors, which are dedicated devices that are capable of the high-speed packet processing that is required for the conversion. Each digital signal processor comprises multiple processing resources, such as processing channels, to handle multiple calls and the different conversion formats across the calls. For example, the conversion formats might be distinguished from one another by codec type, encryption algorithm, payload values, addressing information, or redundancy in the information transmitted. Protocol standards and formulas exist that govern these properties, such as G.711 and G.729 compression/decompression algorithms. Similarly, the media gateway must also perform the conversion in the other direction from a time division multiplexed stream to Internet Protocol packets in a timely manner.\nOne or more media gateways are controlled by a media gateway controller, which provides the call control and signaling functionality for each media gateway and across media gateways. Communication between media gateways and media gateway controllers is achieved by means of protocols such as H.248, Media Gateway Control Protocol (MGCP), and so forth. During a call initialization that involves an Internet Protocol (IP) terminal, the media gateway controller provides to the IP terminal the IP address of the media gateway resource that is handling the call. This enables the IP terminal to specify the proper destination address of the packets that it originates and to recognize the packets that are being sent to the terminal. Alternatively, instead of a call involving an IP terminal, the call could involve another media gateway that exchanges packets with the aforementioned media gateway resource that is handling the call.\nOn occasion, it might become necessary, under certain conditions, to remove an in-service digital signal processor from service. For example, when a technician has to diagnose or physically replace the signal processor, a technician will typically need to “busy out” the processor. What makes this difficult is that the signal processor might be handling the packet streams that correspond to one or more currently active calls between a telecommunications terminal in one type of network, such as the IP network, and a terminal in another type of network, such as the switched telephone network, described earlier. Techniques in the prior art require that the technician issue a graceful busy-out command for a selected signal processor, only to have to wait for its calls to terminate before the processor can be fully removed from service. This can be a lengthy procedure, especially if some calls that are being handled by the selected signal processor are long in duration.\nWhat is needed is a technique to accelerate the removal from service of a signal processor at a media gateway, without some of the disadvantages in the prior art."}
{"text": "The present inventive concept relates to a delay locked loop (DLL), and more particularly, to a DLL for generating a clock signal having a frequency obtained by multiplying the frequency of a reference clock signal, and having an accurate phase delay and a constant duty cycle by compensating for an error occurring due to process variation.\nA DLL may be used to generate an internal clock signal in electronic devices. The DLL usually delays an external clock signal by a predetermined period of time using a delay line and generates an internal clock signal synchronized with the external clock signal. The DLL may be used to generate a timing signal for the operation of a semiconductor memory device, e.g., dynamic random access memory (DRAM), that is, a signal used to read data from the semiconductor memory device or write data thereto.\nTo increase the data transmission speed between electronic devices, it is necessary to increase the frequency of an internal clock signal of the electronic devices. To decrease errors occurring in high-speed data transmission, a data valid window needs to be accurately positioned. For instance, when an internal clock signal having a frequency obtained by multiplying the frequency of an external clock signal is used in a semiconductor memory device, the data transmission speed of the semiconductor memory device can be increased. In addition, errors in high-speed data transmission can be reduced when clock signals having an accurate phase delay and an accurate duty cycle are used.\nWhen the frequency of the internal clock signal increases, however, harmonic locking may occur during the locking operation of the DLL. Also, the influence of process variation is increasing with the microminiaturization of manufacturing technology. Accordingly, it is increasingly difficult to generate a clock signal having an accurate phase delay and a constant duty cycle due to phase delay mismatch between delay lines of the DLL."}
{"text": "The invention relates to a device for providing a breathing gas, in particular to a device for conditioning a breathing gas, for managing the weight of a breathing creature, in particular for stimulating weight reduction. The invention further relates to a method and to a use of a device for conditioning a breathing gas and for providing said breathing gas to a breathing creature, for managing the weight of a breathing creature, in particular for stimulating weight reduction of the breathing creature.\nIn developed countries around the world, there is a tremendous interest in personal fitness, wellness and para medication. At the same time there is a disturbing trend that people, including adults and children are becoming overweight. By overweight, it is meant that the subject has exceeded the acceptable weight range and percent body fat generally considered as healthy determined by factors including, but not necessarily limited to age, height, sex, and body type.\nOverweight produces a wide range of health concerns including sleep apnoea, orthopaedic complications, arterial sclerosis, diabetes, heart disease and also social and psychological problems etc. All these undesirable conditions contribute in development of an unsatisfactory quality of life and in some cases premature death.\nTherefore, many people want to maintain a specific weight or even lose weight to maintain or enhance their physical condition. To lose weight the body needs to burn more energy than is provided by the food intake. Known diet methods include restriction of amounts and kinds of food, exercises and use of diet drugs to reduce or maintain body weight.\nHowever, managing weight via known diet methods takes a long period of time, even when combined with exercises or diet drugs. Many people are not capable of maintaining a reduced food intake and/or an increased exercise level over a long period of time and therefore do not achieve the amount of weight loss aimed for or even gain weight in spite of there efforts. Also, even when the diet results in a loss of weight, for a variety of reasons, most people find it very difficult to maintain significant weight reduction over time.\nTherefore only a few people are able to achieve a significant reduction of weight, and even less are able to maintain this condition."}
{"text": "A variety of batteries employing metallic magnesium anodes are known. Morehouse et. al., in U.S. Pat. No. 2,759,986, taught a magnesium anode cup container lined with a cathode material paste consisting of powdered sulfur, graphite, and magnesium bromide electrolyte. Roche et. al. in Argonne National Lab, ERDA Report No. 76-35, \"Alternative Secondary Cell Systems\", Progress Report, 1976, pp. 53 to 54, discloses a molten electrolyte, magnesium-iron sulfide cell, using FeS.sub.2 cathodes, and MgCl.sub.2 -- NaCl -- KCl electrolyte. The cathode was covered with a stainless steel screen and the anode was covered with graphite fabric. The cell operated at 450.degree. C but short circuited after four cycles due to magnesium dendrite formation.\nOthers in the art have substituted iron sulfide as the cathode material in combination with a lithium anode. These lithium cells provide outstanding high voltage characteristics. Vissers et. al., in U.S. Pat. No. 3,933,521, taught a molten electrolyte, rechargeable, electrochemical cell which required a LiCl -- KCl electrolyte between the anode and cathode and operated at temperatures of about 400.degree. C. The cathode contains iron, cobalt or copper sulfides in a conductive support structure, insulated with a separate boron nitride or yttria cloth separator, and disposed in the center of a housing container. The anode which lines the bottom of the container comprises a 90% porous, compacted, stainless steel or nichel felt, coated with cobalt and impregnated with molten lithium.\nThe use of lithium, due to its high reactivity, requires the use of expensive separator and container materials. In lithium cells, the anode and electrolyte are very expensive. The use of magnesium metal provides severe short circuit complications due to dendritic growth at the magnesium electrode during recharging. Magnesium also reacts chemically with a wide variety of separator materials, particularly oxides, to make them electrically conducting, thus causing shorting of electrodes. What is needed, is an inexpensive, secondary, electrochemical cell, which will still provide good voltage characteristics, and which will not be easily degraded over long time periods."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital telecommunications system, and more particularly, to a DTV receiver and a method of processing a broadcast signal in a DTV receiver.\n2. Discussion of the Related Art\nThe Vestigial Side Band (VSB) transmission mode, which is adopted as the standard for digital broadcasting in North America and the Republic of Korea, is a system that has been developed for the transmission of MPEG video/audio data. However, presently, the technology for processing digital signals is being developed at a vast rate, and, as a larger number of the population uses the Internet, digital electric appliances, computers, and the Internet are being integrated. Therefore, in order to meet with the various requirements of the users, a system that can transmit diverse supplemental information in addition to video/audio data through a digital television channel needs to be developed.\nSome users may assume that supplemental data broadcasting would be applied by using a PC card or a portable device having a simple in-door antenna attached thereto. However, when used indoors, the intensity of the signals may decrease due to a blockage caused by the walls or disturbance caused by approaching or proximate mobile objects. Accordingly, the quality of the received digital signals may be deteriorated due to a ghost effect and noise caused by reflected waves. However, unlike the general video/audio data, when transmitting the supplemental data, the data that is to be transmitted should have a low error ratio. More specifically, in case of the video/audio data, errors that are not perceived or acknowledged through the eyes or ears of the user can be ignored, since they do not cause any or much trouble. Conversely, in case of the supplemental data (e.g., program execution file, stock information, etc.), an error even in a single bit may cause a serious problem. Therefore, a system highly resistant to ghost effects and noise is required to be developed.\nThe supplemental data are generally transmitted by a time-division method through the same channel as the MPEG video/audio data. However, with the advent of digital broadcasting, ATSC VSB digital television receivers that receive only MPEG video/audio data are already supplied to the market. Therefore, the supplemental data that are transmitted through the same channel as the MPEG video/audio data should not influence the conventional ATSC VSB receivers that are provided in the market. In other words, this may be defined as ATSC VSB compatibility, and the supplemental data broadcast system should be compatible with the ATSC VSB system. Herein, the supplemental data may also be referred to as enhanced data. Furthermore, in a poor channel environment, the receiving performance of the conventional ATSC VSB receiving system may be deteriorated. More specifically, resistance to changes in channels and noise is more highly required when using portable and/or mobile receivers."}
{"text": "Costs incurred for producing and circulating publications are often offset by revenue generated by including advertisements within the publication. Advertisements are more successful when properly targeted. Web page advertising has progressed such that advertisements can be targeted to particular viewers. This is not the case for printed publications such as newspapers where all readers of a publication run are presented with the same advertisements. Publishers do not have the tools to dynamically insert different advertising content into different copies of a given publication."}
{"text": "At present, freely deployed wireless access technologies, most notably the IEEE 802.11 family (a.k.a. WiFi), do not support any reliable means to determine the real achievable bit-rate over the link between a wireless terminal (STA) and an access point (AP), or between two terminal stations. This is due to fact that conventional AP beacons broadcast only their bit-rate technical capability, not real estimations of what could be achieved in the particular environment taking into account signal quality and current congestion levels. As a result, the information used by a STA to connect to an AP is inaccurate and incomplete.\nThis limitation in the art poses problems for WiFi users, for many end-user applications require information on the real achievable bit-rate before making a decision of which access point to associate when multiple access points are available. Conventional WiFi technology provides information regarding signal quality to assist in the determination of the estimated bit-rates; however, the conventional bandwidth (bit-rate) estimators, such as Pathcar, require build up of the full connection (i.e., a full association in the WiFi network between AP and STA and an IP connection over the created link) in order to probe by using network packets. As a result, such methods require a substantial amount of latency and overhead and, most importantly, such techniques are not applicable to a determining which AP to build up a connection to in the first place. The future 802.11k extension promises to provide some two-way signal strength information also from the AP to the STA, including history information, but not real bit-rate estimation.\nIn view of the needs in the art for improved bit-rate estimation prior to making WiFi connections between STAs and APs, a simple method is desired to better estimate the download bit-rate between a STA and an AP that takes into account the congestion on the corresponding AP as well as the signal quality. A system and method for achieving congestion load balancing when multiple APs are available, thus ensuring the best download bit-rate to wireless stations, is also desired. The present invention addresses these needs in the art."}
{"text": "The present invention relates to the treatment of dyed textiles, and more particularly to providing denim textiles which have an overdyed and stone washed appearance."}
{"text": "1. Field of the Invention\nThe present invention relates to circuitry for providing electrostatic discharge (ESD) protection. More particularly, the present invention relates to components for providing ESD protection for a circuit with components made using a 2.5 volt semiconductor process having a 5.0 volt compatible CMOS I/O.\n2. Description of the Related Art\nESD can occur when a human body, or machine comes into contact with lead pins of a chip before the chip is placed in a circuit. ESD can also occur after the chip is placed in a circuit and powered up when a charged up lead pin is accidentally connected to ground. A typical circuit used to test a chip to determine if the chip offers adequate protection during an ESD event includes a series connected resistor and capacitor. To simulate the human body, a resistance of approximately 3000.OMEGA. is utilized in the test circuit along with a capacitance of approximately 10 pF. Because the instantaneous voltage from a human body can be as much as 2000 volts, the capacitor is typically charged up to as much as 3000 volts and connected through the resistor to the lead pin of a chip to assure the chip provides adequate ESD protection.\nTo prevent damage to circuitry on a chip during an ESD event, the output buffers are typically designed to go into avalanche breakdown to connect a lead pin receiving an electrostatic discharge directly to a ground pin without affecting other circuitry on the chip. FIG. 1 illustrates a portion of a typical output buffer of a chip which is powered down. The buffer includes an NMOS pull up transistor 102 and an NMOS pull down transistor 104. During power down, the gates of transistors 102 and 104 are connected to a ground pin through resistors. The transistors of the buffer are also connected to a pad 100 which is coupled to a lead pin of a chip.\nWhen an ESD event occurs providing up to 3000 volts to the pad, transistor 104 goes into breakdown and then snap-back to connect the pad to ground to prevent circuit damage. To describe how snap back occurs during an ESD event to provide ESD protection, reference is made to FIGS. 2 and 3. FIG. 2 shows a cross section of a substrate containing the transistor 104 of FIG. 1. FIG. 3 graphs voltage at the pad vs. current through the source to drain of transistor 104.\nIn an ESD event, the voltage at the drain region 200 rises quickly. The inductance 201 will limit how quickly voltage rises on the drain region 200. Initially, with a voltage rising on the drain 200, the gate to drain overlap capacitance 208 begins to pull the gate high. The series resistance 210 of the gate initially protects the gate oxide because the voltage drop across the resistor is subtracted from the voltage that would otherwise appear across the gate oxide. As the gate is pulled up, the channel between the drain 200 and source 202 turns on.\nThe region of highest electric field is the pinch-off region 204 where the conducting electrons reach impact ionization energies, giving rise to electron-hole pair generation. With electroin-hole pair generation, the n+ drain region 200 will further be pulled high, to cause gate aided breakdown by creating a depletion region in the area 204 of the substrate between the gate 206 and region 200. Breakdown occurs initially in region 204 because the highest electric field occurs at this edge of the gate 206. With breakdown initially occurring, electron hole pairs are generated. Holes move away from the pad toward the p- substrate to pull the substrate high. With the p- substrate going high, forward biasing of a diode formed from the p- substrate region to the n+ source region 202 occurs with the n+ region connected to ground through an inductive line 203. With the p- to n+ diode forward biased, electrons will be injected from the region 202 into the p- substrate, which increases the carrier concentration locally. As the carrier concentration increases, the depletion width to maintain a given maximum electric field decreases thus decreasing the voltage drop--hence creating snap back.\nAs current increases through the transistor 104 beyond a primary breakdown point 300, which occurs at snap back, a secondary breakdown can occur at a point 302 as shown in FIG. 3. Secondary breakdown is destructive to the transistor 104 causing it to essentially melt. The size of transistor 104 is typically set to limit current below secondary breakdown when an ESD event occurs. If the size of the single transistor 104 is insufficient to limit current below secondary breakdown when an ESD event occurs, multiple transistors can be utilized as shown in FIG. 4.\nIn FIG. 4, three transistors 401-403 are used to connect the pad to ground. With three transistors and no resistors 411-413, when an ESD event occurs, one of the transistors 401-403 will break down first to limit the voltage potential across the source to drain of remaining transistors in 401-403, so the remaining transistors in 401-403 will not turn on. Current will then increase in the transistor which is turned on to cause secondary breakdown in the transistor turned on. To enable all transistors 401-403 to turn on during an ESD event, ballast resistors 411-413 are included in series with the source to drain path of each transistor 401-403. The ballast resistors have resistance values set so that if one of the transistors 401-403 goes into primary breakdown, the voltage applied across the remaining transistors in 401-403 will be enough to turn on those transistors so that current will not increase in one transistor enough to cause a secondary breakdown.\nFor an ESD event where the pad is pulled negative, referring to FIG. 2, the diode formed by region 200 and the p- substrate will be forward biased causing electron injection into the substrate. A p+ guard region 210 implanted in the substrate near the transistor 104 will complete the circuit to ground and so limits any negative voltage.\nThe structure for a 5.0 volt compatible output buffer made using a 2.5 volt semiconductor process makes providing ESD protection more difficult. A portion of a 5.0 volt compatible output buffer made using a 2.5 volt transistor process is shown in FIG. 5. The circuit includes two pull down transistors 501-502 with source to drain paths connected in series to connect a pad to ground. The circuit also includes two pull up transistors 503 and 504 with source to drain paths connected in series to connect the pad to a pin power supply voltage Vcc. In operation, when a 5.0 volt difference is applied from the pad to ground, or from Vcc to the pad, voltages are applied to the gates of transistors to prevent a voltage potential difference across any of transistors 501-504 from exceeding 2.7 volts. For example, when 5.0 volts is applied from the pad to ground, 3.3 volts is applied to the gate of transistor 502 to cause a 2.5 volt drop from the pad to node n2. With n2 at 2.5 volts, the gate to source, gate to drain, and source to drain voltages of transistors 501 and 502 do not exceed a maximum of 2.7 volts. Circuitry to control the gate voltage of transistors 501-504, as well as additional circuitry for a 5.0 volt compatible I/O buffer made with a 2.5 volt process, is disclosed in U.S. Patent Application Ser. No. 08/912,763, incorporated herein by reference.\nFIG. 6 shows a cross section of a substrate containing either the transistors 503 and 504 of FIG. 5, or transistors 501 and 502 of FIG. 5, with a chip powered down. During an ESD event, when a high voltage is applied to the pad, with the gate of transistor 601 grounded, and the n+ region 616 of transistor 602 grounded, the n+ region 610 will be pulled high to initiate gate aided breakdown. With breakdown beginning, electron hole pairs are separated and the holes will move away from the region 610 toward the p-substrate to pull the substrate high. With the p-substrate going high, forward biasing of a diode formed from the p- substrate region to the n+ source region 616 occurs, since region 616 is grounded. With the p- to n+ (616) diode forward biased, electrons will be injected from the region 616 into the p- substrate to eventually cause snap back.\nHowever, with transistors 501 and 502 of FIG. 5 spaced close enough to allow snap back for ESD protection, snap back can also occur during normal operation. During normal operation, the source of transistor 501 is grounded, a condition needed for snap back. When the output switches from high to low, current through the source to drain of transistor 502 can give rise to hot carrier generation creating electron-hole pairs that can charge up the p- substrate so that snap back can occur.\nTo prevent snap back during normal operation, the transistors 501 and 502 of FIG. 5 can be physically separated in the substrate. However, with transistors 501 and 502 physically separated to an extent needed to prevent snap back during normal operation, snap back will not occur during an ESD event. Current during an ESD event sufficient to cause breakdown at a region 611 in FIG. 6, as well as a region 615 will damage circuitry on a chip containing the buffer.\nIt is, therefore, desirable to provide a new ESD protection mechanism for a 5.0 volt compatible I/O buffer made using a 2.5 volt process, where two series transistors are required from a pad to ground."}
{"text": "The aim of exploration seismology is to obtain an image of the subsurface by probing it with seismic waves at various locations. These waves are generally generated by using airguns in marine, and vibroseis or dynamite in land. They propagate downwards through the subsurface, and are reflected at interfaces between geological layers or refracted within layers. Parts of these waves subsequently propagate upwards to the surface, where they are detected and recorded.\nIn exploration seismology, although the time coordinate is regularly sampled, spatial coordinates are typically irregularly sampled due to the presence of obstacles in land and strong currents in marine. But even for example receivers placed within a marine survey cable or streamer may not be always equidistant. Hence, the inline sampling can be quite irregular.\nThe regularization of seismic data is very important especially in time-lapse survey matching, multiple suppression and imaging. If the irregular nature of the sampling grid is ignored or handled poorly, notable errors are introduced the severity of which may be further amplified at later stages of the seismic processing chain. The problem of signal reconstruction from uniformly spaced data has been investigated in depth. The Whittaker-Kotel'nikov-Shannon sampling theorem states that any signal f(x) can be reconstructed from its uniformly spaced samples if the sampling interval is less than half the period of the highest frequency component in that signal. Thus if f(x) is bandlimited to the wavenumber s/2, which is known as the Nyquist wavenumber, then the sampling theorem provides the following formula to interpolate any function value from uniformly spaced values f(m/s):\n                              f          ⁡                      (            x            )                          =                              ∑                          m              =                              -                ∞                                      ∞                    ⁢                                    f              ⁡                              (                                  m                  /                  σ                                )                                      ⁢            sin            ⁢                                                  ⁢                          c              ⁡                              (                                                      σ                    ⁢                                                                                  ⁢                    x                                    -                  m                                )                                                                        [        1        ]            where sinc(x)=sin(πx)/πx. Thus, when the sampling rate is sufficient and there is no aliasing, the sampling theorem provides a way to reconstruct the signal “exactly” from its uniformly spaced samples. To satisfy requirements of the sampling theorem, the signal should be sampled at a rate greater than twice the Nyquist rate, i.e., s. The seismic signal may be given by:\n                                          f            ⁡                          (              x              )                                ≈                                    f              L                        ⁡                          (              x              )                                      =                  σ          ⁢                                    ∑                              m                =                0                                            L                -                1                                      ⁢                          Δ              ⁢                                                          ⁢                              x                m                            ⁢                              f                ⁡                                  (                                      x                    m                                    )                                            ⁢              sin              ⁢                                                          ⁢                              c                ⁡                                  (                                      σ                    ⁡                                          (                                              x                        -                                                  x                          m                                                                    )                                                        )                                                                                        [        2        ]            where Δxm is the Jacobian weight, i.e., Δxm=xm+1−xm and f(xm) the value of the seismic data at irregular offset xm. It is important to note that, when Δxm=1/σ the sinc interpolator of eq. [2] is exact since\n                              f          ⁡                      (            x            )                          =                              ∑                          m              =                              -                ∞                                      ∞                    ⁢                                    f              ⁡                              (                                  m                  /                  σ                                )                                      ⁢            sin            ⁢                                                  ⁢                          c              ⁡                              (                                  σ                  ⁡                                      (                                          x                      -                      τ                                        )                                                  )                                                                        [        3        ]            by Whittaker-Kotel'nikov-Shannon theorem. On the other hand, when Δxm is not equal 1/σ, the sinc interpolator provides only a crude approximation to the continuous signal.\nFor such cases, it has been found that a better approach is to invert equation [1] for the desired uniformly spaced signal values f(m/σ). This inversion when written in matrix notation yields\n                    h        ≡                                                                                        f                  ⁡                                      (                                          x                      1                                        )                                                                                                                        f                  ⁡                                      (                                          x                      2                                        )                                                                                                      M                                                                                      f                  ⁡                                      (                                          x                      L                                        )                                                                                                ≈                                                                                                              s                    11                                                                                        s                    11                                                                    Λ                                                                      s                                          1                      ⁢                                                                                          ⁢                      L                                                                                                                                        s                    21                                                                                        s                    11                                                                    Λ                                                                      s                                          2                      ⁢                                                                                          ⁢                      L                                                                                                                    M                                                  M                                                  O                                                  M                                                                                                  s                    11                                                                                        s                    11                                                                    Λ                                                                      s                    LL                                                                                            ·                                                                                                      f                    ⁡                                          (                      0                      )                                                                                                                                        f                    ⁡                                          (                                              1                        /                        σ                                            )                                                                                                                                                                                                                                                            f                    ⁡                                          (                                                                        (                                                      L                            -                            1                                                    )                                                /                        σ                                            )                                                                                                                      ≡        Sg                            [        4        ]            where σ/2 is the bandwidth of the signal f(x) and S is the sinc matrix with entries sij=sinc(σ(xi−j/σ). If the matrix S is well conditioned than the seismic data at regular offsets can be computed by standard matrix inversion:g=S−1h  [5]Otherwise, a least squares minimum norm inversion can be used:g=(STW1S+W2)−1STW1h  [6]where W1 is usually chosen as a diagonal matrix whose mth diagonal entry is the Jacobian weight Δxm=xm+1−xm and W2 is usually chosen as a small multiple of identity matrix, i.e., W2=ε2I.\nThis necessitates the use of robust and efficient techniques for seismic data regularization.\nAn interpolator in accordance with equation [6] has been described in: Yen J. L., 1956, On nonuniform sampling of bandwidth-limited signals, IRE Trans. Circuit Theory, CT-3, 251-257 (1956). It is therefore sometimes referred to as Yen's Interpolator of Type 1.\nMany interpolators used in seismics are variations of the Yen's interpolator for example those described in: Duijndam, A. J. W., Schonewille, M. A., Hindriks, C. O. H., Irregular and sparse sampling in exploration seismology (chapter of L. Zhang: Nonuniform sampling: theory and practice, Kluwer Academic/Plenum Publishers, New York, USA (2001)). The regularization in Duijndam is formulated as a spectral domain problem. There it is proposes to estimate the spectrum of the signal by using non-uniform Fourier transform of the irregular samples, and then reconstruct the regular samples by using inverse discrete Fourier transformation. It can be shown that, this regularization technique is exactly equivalent to Yen's Interpolator of Type 1.\nA variant of this method is described in: Zwartjes, P. M. and M. D. Sacchi, 2004, Fourier reconstruction of non-uniformly sampled, aliased data: 74th Ann. Internat. Mtg.: Soc. of Expl. Geophys., 1997-2000. However in Zwartjes there is proposed least squares inversion of the Fourier transformation instead of using an inverse discrete Fourier transform. To this purposes a cost function is defined, which also involves a non-quadratic penalty term to obtain a parsimonious model.\nAnother example is: Hale, I. D., 1980, Resampling irregularly sampled data, Stanford Exploration Project, SEP-25, 39-58. Hale's method is based on the more general version of the Yen's Interpolator, where a space limited signal assumption is not used. In that case, the uniform samples f(m/σ) can still be computed by solving a matrix equation similar to [6]. What Hale suggests is to replace the entries in the inverse matrix by their locally computed approximations.\nThe interpolators based on Yen's 1st theorem usually provide satisfactory results on non-aliased signals with little high-wavenumber content. However their performance degrades significantly when the interpolated signal has a substantial amount of high wavenumber spectral content. Another shortcoming of the interpolators based on Yen's 1st theorem is that in order to solve [4], at least as many irregular sampling positions as regular sampling positions are required. Hence, if some seismic traces are dropped out, traces which reside at further locations must be used to solve the system of equations given by [4]. Usually this degrades the accuracy of the interpolated sample values.\nFurther, although Yen's 1st interpolator is exact for infinite length signals, it is an approximation when only a finite extent of the signal is available for interpolation.\nGiven the problems of the existing interpolators it remains an object to find improved interpolates capable of interpolating data received by receivers at irregular locations to regulars sampling locations."}
{"text": "1. Field of the Invention\nThe present invention relates an antenna, and, more particularly, to a triple-band antenna for different frequency bands, which is designed for an increased low frequency bandwidth and different intermediate frequency bandwidths.\n2. Description of the Related Art\nWith the rapid growth of wireless communication technologies, standard signal frequency antennas are now insufficient, and so multiple frequency antennas have become the technology of choice. A multiple frequency antenna is usually used in a portable electronic device that supports wireless communication functions, such as a notebook, a mobile phone or a PDA. Since these electronic devices are all very thin and light, it is necessary to have small-volume multiple frequency antennas. However, usually when the antenna has a smaller volume, its reception efficiency is also reduced, and multiple frequency antennas may have narrow frequency bandwidths at different frequency locations. Therefore, the design needs to compromise between volume and reception efficiency. Moreover, the standard multiple frequency antenna with an intermediate frequency band reception ability may also fail to have a broadband response due to the design.\nIt is therefore desirable to provide a triple-band antenna to mitigate and/or obviate the aforementioned problems."}
{"text": "Ambulances are able to carry only a very small number of patients (generally only one or two), making it extremely complex to transport all the victims or patients that have treatment or evacuation needs when a mass casualty event or disaster occurs. Currently, when such an event occurs, many ambulances are called to transport the victims, and, given the limited number of ambulances generally available within a given geographic area, the total number of patients that can be transported is also greatly limited. Furthermore, with so many victims needing transportation from one mass casualty event, if all conventional ambulances are used for transport, people in other coverage areas might suffer because ambulances are not available to handle other emergencies.\nA similar problem is encountered when an evacuation of a large number of patients is required. For example, the ability to relocate the patients from a hospital or a nursing home that may be in the path of an approaching hurricane also would be limited by the availability and number of ambulances available.\nSimilarly, if there is a need to treat a large number of victims at a given location, the only option available now is to transport the victims to a remote site for treatment, and movement of the large number of victims is similarly impacted by the limited number of ambulances. Alternatively, a tent-like treatment facility might be set up, but the suitability of such a system would depend on the disaster conditions (e.g., a tent would not hold up well during windy conditions; climate extremes might be difficult to control; once patients are stabilized in the tent-like structure, they would still need to be transported to a hospital or other facility for extended care).\nOther ambulance-type vehicles, such as those used in the military, use J-type hooks or loop-type straps for holding stretchers, and loading a stretcher onto such a J-hook type holder or loop-type straps requires up to six people to maneuver the stretcher into storage position.\nFurthermore, if there is a need to transport a large number of casualties, multiple ambulances or hearses would be needed."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a system managing method and apparatus and, more particularly, to a system managing method and apparatus using both a managed object interface and a non-managed object interface.\nIn a system such as an exchange system supporting a plurality of kinds of interfaces including a managed object (MO) interface and a non-managed object (non-MO) interface, it is desired to simplify a process using the MO interface and the non-MO interface. The MO interface includes common management information protocol (CMIP). A TL1 command and a graphical user interface (GUI) are included in the non-MO interface.\n2. Description of the Related Art\nFIG. 1 is an illustration for explaining a conventional system management model. The system management model shown in FIG. 1 is an open system interconnection (OSI) system which enables interconnection between different computer network systems. In the system management model shown in FIG. 1, communication between a managing system 41 (manager 43) and a managed system 42 is specified. The managed system 42 includes a plurality of managed objects (MOs) 45 and an agent 44 which operates the managed objects 45.\nThe managing system 41 (manager 43) may correspond to a managing apparatus such as a maintenance terminal apparatus provided with processors and memories. The managed system 42 may correspond to a switching apparatus provided with processors and memories. In the system management model shown in FIG. 1, resources to be managed are abstracted so as to manage the resources as the managed objects. The managed objects include concepts of a managed object class (MO class) and a managed object instance (MO instance). The MO class defines a characteristic of a managed object. The MO class defines attributes, attribute groups, notifications, operations and behavior. The MO instance represents a specific object defined by the MO class.\nFor example, the above-mentioned CMIP is defined as a protocol for transmission and reception of management commands and notification information of the managed objects 45 between the managing system 41 (manager 43) and the agent 44 of the managed system 42. The CMIP is a protocol for use on an individual managed object basis. A simple network management protocol (SNMP) is well-known as the CMIP. An interface between the manager 43 and the agent 44 is referred to as the MO interface. That is, the MO interface specifies communication between the managing system 41 and the managed objects 45 via the agent 44 such as transmission of managing operation commands from the managing system 41 (manager 43) to the managed objects 45 or transmission of notification from the managed objects 45 to the managing system 41.\nSuch a specification of communication is known as a common management information service (CMIS). The managing operations according to the CMIS includes, for example, M-GET which is acquisition of an attribute value, M-SET which is setting of an attribute value, M-ACTION which is an operation on the managed object, M-CREATE which is creation of the managed object, M-DELETE which is deletion of a managed object and M-CANCEL-GET which is cancellation of the M-GET. Notification according to the CMIS includes M-ENENT-REPORT which is a notification sent from a managed object.\nOn the other hand, the TL1 command or the GUI does not handle resources as the managed object. Such an interface is referred to as the non-MO interface. The non-MO interface is used in various systems including a switching system. In the TL1 command, invoked process and parameters for the process are specified.\nFIG. 2 is an illustration for explaining another conventional system management model in which both the MO interface and the non-MO interface are used. The management system model shown in FIG. 2 includes a managing system (manager) 51 for the MO interface, a managed system 52, an agent 54, a common processing unit 55 and a managing system (manager) 56 for the non-MO interface.\nThe managed system 52 is provided with the common processing unit 55 so that a managing operation for both the MO interface managing system 51 and the non-MO interface managing system 52 can be performed. The common processing unit 55 has a structure based on either the MO interface or the non-MO interface.\nIf the common processing unit 55 has a structure suitable for the MO interface in the above-mentioned system in which both the MO interface and the non-MO interface are supported, the unit of process requested of the common processing unit 55 is the MO class and a full distinguished name (FDN). The FDN is a symbol used for designating the MO instance which belongs to one of the MO classes in the MO interface. That is, the FDN is an identification number of the MO instance managed as a hierarchical structure according to a hierarchical relationship between a plurality of MO classes.\nIn this case, a function of the common processing unit 55, which function is in accordance with the CMIP command including the MO class name and the FDN, can be correlated to a command sent from the MO interface. However, an appropriate function of the common processing unit 55 cannot be provided in response to a command sent from the non-MO interface since the command does not include the MO class name and the FDN. In such a case, a process is required for relating the non-MO command to a function of the common processing unit 55 that is based on the MO interface. Thus, the agent 54 performs such a process.\nOn the other hand, if the common processing unit 55 has a structure suitable for the non-MO interface, the function of the common processing unit 55 is constituted by the unit of command handled by the non-MO interface. Thereby, an appropriate function of the common processing unit 55 can be provided in response to a request from the non-MO interface. However, a command in the MO unit from the MO interface must be processed so as to be related to a command unit in the non-MO interface. In such a case, since the command unit of the non-MO interface has a course tendency as compared to that of the MO interface, there is a problem in that an appropriate function of the common processing unit 55 cannot be provided in response to the request from the MO interface.\nAdditionally, if a pointer attribute defined by the MO interface is handled in the common processing unit 55, it is considered to set a definition of various pointers that indicates a relationship between the instances so as to perform processes according to the definition. However, in order to respond to the request, a large process load is applied to the common processing unit 55. Additionally, there is a problem in that a long time and labor are required to develop a program performing such an operation.\nIt is a general object of the present invention to provide an improved and useful system managing method and apparatus in which the above-mentioned problems are eliminated.\nA more specific object of the present invention is to provide a system managing method and apparatus in which a load applied to a common processing unit is decreased so as to improve a managing ability of the system.\nIn order to achieve the above-mentioned objects, there is provided according to the present invention a system managing method and apparatus for managing a system having a plurality of interfaces including a managed object interface and a non-managed object interface, the managed object interface handling resources as managed objects, the non-managed object interface handling resources as non-managed objects. A common processing unit performs operations common to the plurality of interfaces. A managed object interface agent is connected to the managed object interface so as to manage pointer attributes independently from the common processing unit. A non-managed object interface agent is connected to the non-managed object interface. The managed object interface agent includes a knowledge database and an agent database, the knowledge database storing information regarding a pointing instance and a pointed instance defined by each pointer attribute, the agent database holding the resources as instances of the managed object type so as to delete, add or update a value of each pointer attribute according to a process request.\nAccording to the present invention, the pointer attributes can be processed solely by the managed object interface agent. That is, when a process request is transferred from the non-managed object interface agent to the common processing unit, the common processing unit determines whether or not the request is for the resources managed by the managed object interface. If the request is for the resources managed by the managed object interface, the common processing unit transfers the process request to the managed object interface agent so that the request is processed by the managed object interface agent. Thus, the common processing unit is not required to process the request for the resources managed by the managed object interface agent. That is, the common processing unit is required to perform only processes related to the non-managed object and basic processes common to the managed object and the non-managed object. This reduces a load to the common processing unit when a plurality of interfaces including the managed object interface and the non-managed object interface are present."}
{"text": "The present disclosure pertains to thermostats and particularly to mode changes between heat and cool of thermostats."}
{"text": "1. Field of the Invention\nThe present invention relates to image processing technology employing shooting information at the time of shooting, and image processing control information that designates image processing conditions to be used when performing image processing on image data.\n2. Description of Related Art\nTo date, there have been proposed techniques for performing image processing of image data created by an image data generating apparatus such as a digital still camera (DSC), which involve performing image processing with reference to shooting information (shooting conditions) at the time of shooting, and image processing control information that designates processing conditions to be used during image processing. Shooting information and image processing control information is typically described in the header of the image data, and is associated with a particular set of image data.\nAccording to this image processing technique, where shooting information is used, it becomes possible to perform on image data image processing that is adapted to the shooting conditions, thus preventing image processing that is inappropriate for a particular type of photographed scene, for example. Where image processing control information is used, it becomes possible to perform on image data image processing based on certain image processing conditions contemplated at the time of shooting, so that output images reflect the photographer's intention or the DSC vendor's design.\nWith the prior art image processing apparatus, however, since the idea was to use either shooting information or image processing control information, if an image processing apparatus were designed to be capable of image processing using image processing control information, it could not utilize shooting information, even if such information were associated with the image data.\nAlso, in the event that a portion of one set of information cannot be used, if the other set of information includes information corresponding to the information that cannot be used, it will not be possible to use the aforementioned information. Accordingly, there is a need to perform image processing in an appropriate manner, that mutually utilizes, in an effective manner, shooting information and image processing control information."}
{"text": "This disclosure generally relates to systems and methods for wireless communications, and more particularly relates to systems and methods for connecting devices through intermediate nodes.\nIn various situations, stable and reliable cellular, satellite and Internet connections can all be unavailable. In order to provide connectivity to wireless devices without cellular, satellite and Internet connections, wireless devices may communicate between each other through a mesh of devices connected through a local mesh network.\nA wireless communication network can be realized through the use of a mesh network. A mesh network is a distributed network having multiple wireless mesh points or wireless nodes. Each mesh wireless node in the mesh network may act as a relay node capable of receiving wireless signals and relaying the wireless signals to the next mesh wireless node in the mesh network. A wireless signal may proceed from an origination mesh wireless node to a destination mesh wireless node by hopping from one mesh wireless node to another mesh wireless node. Existing systems and methods for routing wireless signals through the mesh network need to insure that the signals are routed efficiently from their origination mesh wireless node to their destination mesh wireless node."}
{"text": "As biotechnology advances, the making of fermentation products by use of bioreactors utilizing immobilized microorganisms is under study in the fields of brewing for malt alcohol beverages (beers), wines, sake, vinegar, soy sauce, and the like. When a bioreactor is used as such, the following are expected:\n1) Since a high concentration of yeast is immobilized to carry out fermentation, the brewing is completed rapidly, so that the brewing period can be shortened, whereby the number of manufacturing tanks and the cost of construction can be lowered.\n2) Since continuous fermentation is possible, it is unnecessary to charge and collect yeast.\nConventionally, when a beer or the like is manufactured by use of a bioreactor, however, the amount of amino acids and diacetyl (DA) has become greater, and the amount of ester has become smaller in the product as compared with a product manufactured by a traditional batch fermentation method. As a result, the product obtained by use of the bioreactor has been problematic in that it is disadvantageous in flavor, and this problem has been remarkable in particular when the bioreactor is used for primary fermentation (main fermentation).\nFor overcoming such a problem, studies have conventionally been carried out. Japanese Patent Application Laid-Open Gazette No. HEI 7-123969 discloses a method in which a fermentation liquid is circulated in a continuous fermentation method using a fluidized bed type reactor. Though the consumption of amino acids in the fermentation process is ameliorated by this method, the amount of diacetyl which causes raw odor or immature odor is not lowered sufficiently, whereby the resulting product still has room for improvement in terms of flavor. Also, manufacturing methods in bioreactors including this method have been problematic in that the number of floating yeast cells upon the end (completion) of primary fermentation is small in general, the fermentation rate in the process of fermentation is unstable, and they are hard to control.\nMeanwhile, it has conventionally been known that yeast for making a fermentation product such as beer is required to have flocculation ability (agglutinability) to a certain extent. Here, flocculation ability refers to a property in which yeast cells flocculate as a mass upon the end of fermentation and are separated from the liquid so as to float up or sediment. If the flocculation ability is too high, then the flocculation occurs at an earlier stage, so that the contact between the fermentation liquid and yeast is broken, whereby fermentation becomes insufficient. If the flocculation ability is too low, by contrast, then yeast floats in the liquid for a long time, so that the amount of recovery of yeast decreases, which takes unnecessary time and labor for separating and collecting yeast by centrifuge or the like, for example. Therefore, it has conventionally been common knowledge to use yeast having a certain degree of flocculation ability (flocculent yeast) for making a fermentation product. Such flocculent yeast (agglutinative yeast) is used in the above-mentioned method disclosed in Japanese Patent Application Laid-Open Gazette No. HEI 7-123969 as well.\nIn view of the above-mentioned problems in the prior art, it is an object of the present invention to provide a method in which, when making a fermentation product by use of a bioreactor utilizing an immobilized microorganism, the fermentation rate in the fermentation process can be held constant, and the number of floating yeast cells upon the end of fermentation can stably be maintained at a level higher than the conventional one, whereby, when making a malt alcohol beverage by use of a bioreactor in particular, the amount of diacetyl in the fermentation liquid and final product can be lowered sufficiently and so forth, so as to improve the flavor of final product."}
{"text": "1. Field of the Invention\nThe present invention relates to failure prediction for disk drives. In particular, the present invention relates to predicting failure of remote disk drives at a central processing facility using an evolving disk drive failure prediction algorithm.\n2. Description of the Prior Art\nPredicting impending failure of disk drives deployed in-the-field helps protect against catastrophic data loss, for example, by suggesting an end user back-up and/or replace the failing disk drive. Another application of disk drive failure prediction is preemptive maintenance of consumer electronics devices that utilize disk drives, such as personal digital assistants (PDAs), digital cameras, personal video recorders (PVRs), and the like. For example, if the failure prediction algorithm indicates impending failure, the end user may return the consumer electronic device to the vendor for repair or replacement.\nIn the past, failure prediction algorithms have been implemented within each disk drive wherein the disk drive itself issues a warning when failure is imminent. An example of a failure prediction algorithm implemented within a disk drive is the Self-Monitoring Analysis and Reporting Technology (SMART), which monitors a number of quality metric values (e.g., fly-height), and reports an impending failure if any one of the quality metric values exceeds some predetermined threshold. A problem with internally implemented failure prediction algorithms, such as SMART, is their limited degree of sophistication which increases the likelihood that impending failures are not detected leading to catastrophic data loss, or that impending failures are falsely detected (false alarm) leading to unnecessary and costly remedial action.\nThere is, therefore, a need to improve upon disk drive failure prediction so as to increase the failure prediction accuracy while minimizing the number of false alarms."}
{"text": "Field of the Invention\nThe present invention relates, in general, to toys. A long-time favorite toy is formed with two buttons or disks rotatably mounted on a continuous loop of string. The ends of the loop of string are grasped by the hands of the user and the string is rotated about a horizontal axis between the user's hands to wind up the string. Then, in and out reciprocating movement of the ends of the strings causes the string to unwind and rewind which forces the disks to oscillate back and forth along the string and strike each other.\nWhile this toy provides long periods of enjoyment, it is not without a disadvantage. Since the user employs his or her fingers to grasp the ends of the loop of string, the fingers are susceptible to being pinched by the winding and unwinding action of the string. Looping the ends of the string through rings which the user grasps eliminates any pinching of the user's fingers. However, the string is typically tied to one of the rings which causes unequal tension in the upper and lower portions of the string during the unwinding and rewinding of the string and prevents the string from continuously rewinding. Thus, long term repeated use of such a toy has not been possible.\nThus, it would be desirable to provide a rotatable disk string toy which includes means to facilitate easy handling of the string. It would also be desirable to provide a rotatable disk string toy with string holding means which is easily attached to the toy. It would also be desirable to provide a rotating disk string toy which includes string holding means having an inexpensive cost. It would also be desirable to provide a rotating disk string toy which is durable for a long useful life. Finally, it would be desirable to provide a rotatable disk string toy which provides a unique visual appearance when in use."}
{"text": "For aluminium production, technology based on a carbothermic process is promising and offers the prospect of an alternative to the Hall-Héroult electrolytic technology. A successful carbothermic process would have the potential to reduce capital investment requirements by 50 to 70% and operating costs by 25 to 35% compared to the electrolytic route. Also, the problem of fluoride emission would be obviated, while the quantity of generated carbon containing gases would be substantially lower than for electrolytic production of aluminium.\nAttempts to produce aluminium by a carbothermic process have been made for in excess of 100 years. However, optimisation of a carbothermic process to enable successful commercial production of aluminium is yet to be achieved.\nProcesses investigated to this stage require temperatures in excess of 2000° C. and accurate control of reactants and products at different complex stages. The stages include:                (a) reaction of alumina and carbon to produce aluminium carbide at above 2000° C.;        (b) reaction of the aluminium carbide with alumina to produce aluminium metal at above 2150° C.; and        (c) separation of the aluminium from remaining materials.        \nChallenges to be met in such carbothermic process include successfully recovering the high level of volatilized aluminium, reducing the level of refractory loss, the difficulties of transferring materials between stages and the problem of generation of a high volume of carbon monoxide. Such issues are inevitable at operating temperatures as high as 2000 to 2250° C.\nReactions central to the carbothermic processes are:2Al2O3+9C→Al4C3+6CO,  (1) andAl2O3+Al4C3→6Al+3CO  (2)\nThese reactions give the overall reaction of:Al2O3+3C→2Al+3CO  (3)\nEarlier work on the production of aluminium by these reactions is illustrated by U.S. Pat. Nos. 1,219,797 and 1,222,593 both to Barnet et al; U.S. Pat. Nos. 2,090,451 and 2,255,549 both to Kruh; U.S. Pat. No. 2,755,178 to Rasmussen; U.S. Pat. No. 2,776,884 to Grunert alone; and U.S. Pat. No. 2,829,961 to Miller at al; and U.S. Pat. No. 2,974,032 to Grunert. More recent work has been directed to reacting alumina and carbon in a molten bath having a molten slag of aluminium carbide and alumina. The molten bath usually operates with two zones, in a first of which aluminium carbide is generated, and a second to which the carbide passes to be reacted with alumina to produce metallic aluminium. This work is illustrated by U.S. Pat. No. 4,385,930 to Persson; U.S. Pat. No. 6,440,193 to Johansen et al; U.S. Pat. No. 6,475,260 to LaCarmera; U.S. Pat. No. 6,530,970 to Lindstad; U.S. Pat. No. 6,849,101 to Fruehan et al; and US patent application publication 2006/0042413. Also of interest are the publications: “Carbothermal Production of Aluminium” by Motzfeldt et al, published in 1989 by Aluminium-Verlag GmbH of Düsseldorf, Germany; and “Aluminium Carbothermic Technology” submitted to US Department of Energy under Cooperative Agreement Number DE-FC36-00ID13900 by M J Bruno and Alcoa Inc, and dated 31 Dec. 2004."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device having circuits structured by electric field effect transistors (FETs), for example, thin film transistors (TFTs), and to a method of manufacturing the semiconductor device. The present invention relates, for example, to a semiconductor device, typically a liquid crystal display panel, and to an electronic device in which such a semiconductor device is mounted as its component.\nNote that, throughout this specification, the term electro-optical device indicates general devices for performing shading display by changing an electrical signal, and that liquid crystal display devices and display devices using electroluminescence (EL) are included in the category of electro-optical device.\nNote also that, throughout this specification, the term element substrate indicates general substrates on which active elements such as TFTs and MIMs are formed.\n2. Description of the Related Art\nTechniques of structuring thin film transistors using semiconductor thin films (having thicknesses on the order of several nm to several hundred nm) formed on a substrate having an insulating surface have been focused upon in recent years. The thin film transistors are being widely applied to electronic devices like ICs and semiconductor devices, and in particular, their development has accelerated rapidly as switching elements of liquid crystal display devices.\nLiquid crystal display devices are known to be roughly divided into active matrix types and passive matrix types.\nA high-grade image can be obtained with active matrix liquid crystal display devices using TFTs as switching elements. Active matrix applications are generally to notebook type personal computers, but they are also expected to be used in televisions for a home and in portable information terminals.\nThe active matrix liquid crystal display devices are generally driven by line inversion drive. With the line inversion drive, for example source line inversion drive, the polarity of voltages applied to adjacent source lines differs, as shown in FIGS. 37A and 37B, and the polarity of the voltage applied to each source line changes each frame. FIGS. 37A and 37B show the polarity of voltages applied to pixels during the source line inversion drive. Drive in which the polarity of the voltage differs for each adjacent source line is referred to as the source line inversion drive. Drive in which the polarity of the voltage differs for each adjacent gate line is referred to as gate line inversion drive.\nFIG. 10 shows schematically a cross section of a pixel portion of a liquid crystal display device. An electric field formed between pixel electrodes 102a and 102b formed on a substrate 101, and an opposing electrode 103 formed on an opposing substrate 104, as shown in FIG. 10, is referred to as a vertical direction electric field 105 in this specification. Further, an electric field formed between the adjacent pixel electrodes 102a and 102b is referred to as a horizontal direction electric field 106 in this specification.\nLiquid crystals in the vicinity of the pixel electrodes orient themselves along the horizontal direction electric field if the line inversion drive is performed, the liquid crystal orientation in edge portions of the pixel electrodes becomes nonuniform, and disclinations develop. In order to obtain a good quality black level, light shielding films for covering the disclinations are necessary. However, the aperture ratio drops if the disclinations are covered by the light shielding films. It is necessary to come up with a scheme in which a good quality black level can be obtained, and as little disclination as possible develops when displaying a high aperture ratio, bright image. Note that, in this specification, liquid crystal orientation irregularities developing due to differences in the direction of the pre-tilt angle, and differences in the twist direction, at liquid crystal orientation film interfaces are referred to as xe2x80x9cdisclinationsxe2x80x9d. Further, regions having different brightnesses produced due to an irregular orientation state of the liquid crystals when a polarization plate is formed is referred to as xe2x80x9clight leakagexe2x80x9d.\nIn particular, the occupied ratio of pixels in which disclinations and light leakage developing due to horizontal direction electric fields is large enough that it cannot be ignored in liquid crystal display devices in which the pixels are formed at a very fine pitch, such as that of a projecting liquid crystal display device. Further, these disclinations and light leakage are expanded when projected onto a screen with the projecting liquid crystal display device, and therefore whether or not the light leakage and disclinations can be suppressed is vital in maintaining contrast.\nAn object of the present invention is to provide an element structure such that liquid crystal disclination and light leakage can be stopped in an active matrix liquid crystal display device.\nThe following measures are taken in order to solve the above stated problems with the conventional technique.\nOverlapping edge portions of pixel electrodes with predetermined height convex portions.\nFIG. 2 is a cross sectional diagram of a simulation model. The present invention utilizes moving disclinations and light leakage, which are caused when a voltage is applied to liquid crystal 202, to edge portions of pixel electrodes by forming edge portions of pixel electrodes 203a and 203b on a first substrate (not shown) so as to overlap with convex portions 204 formed on a level surface as shown in FIG. 2. An opposing electrode 201 is formed in an opposing substrate 207.\nNote that, in this specification, the convex portions are formed selectively below the pixel electrodes. Regions in which the pixel electrodes overlap with upper edge portions of the convex portions are referred to as first regions (a) of the pixel electrodes. Regions in which the pixel electrodes are formed in side portions of the convex portions are referred to as second regions (b) of the pixel electrodes. Regions in which the pixel electrodes are formed on a level surface, and which contact the second regions of the pixel electrodes, are referred to as third regions (c) of the pixel electrodes.\nFurther, the height (h) of the convex portions is the maximum value of the length of a vertical line formed from the upper edge portions of the convex portions to the level surface on which the convex portions are formed.\nIn addition, a cell gap (d) is the distance from the opposing electrode formed on the opposing substrate (second substrate) to the third region of the pixel electrode.\nAn inter-pixel electrode distance (s) is the distance between the first regions of mutually adjacent pixel electrodes.\nConventionally, if there are convex portions in the orientation surface of liquid crystals, then the orientation of the liquid crystal is disordered and light leakage develop at the convex portions, and therefore it is thought that a liquid crystal orientation surface should be as level as possible. However, the applicants of the present invention found that when the first regions of the pixel electrodes formed on the convex portions having a predetermined height, and the second regions of the pixel electrodes formed on the side portions of the convex portions having a predetermined height, are present, simulation results on the liquid crystal orientation show that liquid crystal orientation irregularities caused by the horizontal direction electric field in driving a liquid crystal display device can be reduced. Specifically, locations at which disclinations and light-leakage appear are in the edges of the pixel electrodes during black display.\nThis phenomenon is explained by the schematic diagrams of FIGS. 1A to 1C which show the principles of the present invention. The liquid crystal orientation method is taken as a TN method. FIGS. 1A to 1C show liquid crystal orientations when driving a liquid crystal display device by a 5 V or xe2x88x925 V video voltage with line inversion drive. An orientation film is not shown in the figures.\nFirst, as shown in FIG. 1A, if edge portions of the pixel electrodes 203a and 203b formed on the first substrate (not shown in the figures) are formed on the convex portion 204, then it becomes more difficult for disclinations to develop, compared to the case where no convex portion presents. Nevertheless, if the height of the convex portion is low, then the influence of a horizontal direction electric field formed between the pixel electrode 203a and the pixel electrode 203b is strong with respect to a vertical direction electric field formed between the pixel electrodes 203a and 203b and the opposing electrode 201 formed on the second substrate (not shown in the figures). Liquid crystal molecules 208 in the vicinity of the convex portion orient with an inclination in a diagonal direction with respect to the substrate surface. Thus, light leakage can be seen below a cross Nicol polarization plate. Further, a pre-tilt angle of liquid crystal molecules 209 in the vicinity of the interface of the orientation film is determined by rubbing directions 205 and 206, and therefore a disclination 210 having high light strength is formed at the point where the direction of the horizontal direction electric field in the vicinity of the interface differs from the direction of the liquid crystal molecules in the vicinity of the interface due to rubbing.\nHowever, the higher the convex portion becomes, the more the position at which the disclination 210 seen in FIG. 1A appears changes to the edge of the pixel electrode. In addition, if the convex portion 204 overlapping with the edge portions of the pixel electrodes 203a and 203b is made higher, then a vertical direction electric field formed between a first region 215 of the pixel electrode and the opposing electrode 201 becomes stronger with increasing the height of the convex portion, as shown in FIG. 1B, and the influence of the horizontal direction electric field becomes weaker. Several vertical direction electric fields are formed almost in the substrate surface by an electric field formed between a second region 212 of the pixel electrode, among the pixel electrodes 203a and 203b, formed in the side portion of the convex portion and the opposing electrode 201. TN method liquid crystals are positive type liquid crystals, and therefore the longitudinal axes of liquid crystal molecules 211 orient parallel to the electric fields. Disclination and light leakage in the vicinity of the convex portion is thus reduced.\nNext, as shown in FIG. 1C, if the height of the convex portion 204 is made higher, then the liquid crystals on the pixel electrode 203b form an electric field having a direction diagonal with respect to the substrate surface, the same direction as the pretilt angle, between the opposing electrode and the second region of the pixels. Thus, liquid crystal molecules 213 orient diagonally with respect to the substrate surface, along the electric field, and a light leakage having a width which cannot be ignored is formed in the vicinity of the convex portion. On the pixel electrode 203a, an electric field having a direction diagonal with respect to the substrate surface and formed between a second region 216 of the pixel electrode and the opposing electrode is opposite to the direction of the pre-tilt angle of the liquid crystals, and therefore it is difficult for liquid crystals 217 to follow the electric field, and formation of light leakage becomes relatively difficult. However, the width of the light leakage formed due to the electric field formed between the opposing electrode and the second region 212 of the pixel electrode 203b becomes wider, and overall, the aperture ratio drops.\nIt can be understood from the above that the optimal value of the height of the convex portion exists in order to reduce both the disclination of black display and the width of the light leakage when the edge portions of the pixel electrodes are formed on the convex portion. If the convex portion is too tall, then overall, the width of the light leakage will become larger (see FIG. 1C.) The structures of FIG. 1A and FIG. 1B can increase the aperture ration. The subsequent simulation results substantiate these principles.\nThe optimal value of the height of the convex portion can be considered to be determined with a cell gap (namely, an element for determining the strength of the vertical direction electric field) as a parameter.\nThe applicants of the present invention performed simulations, and confirmed the optimal value of the convex portion height.\nDisclination and light leakage due to horizontal direction electric fields become problems in particular when the pixel area is small, and the proportion of the pixel occupied by the disclination and the light leakage is large enough that it cannot be ignored. In other words, mainly for cases where the device is used as a projection liquid crystal display device. Projection liquid crystal display devices have a small pixel pitch, and inevitably the distance between the pixel electrodes is small at 4.0 xcexcm or less. The applicants of the present invention performed simulations focusing on inter-pixel electrode distances of equal to or less than 4.0 xcexcm in order to reduce the disclination and the light leakage in the projection liquid crystal display device.\nA simulation model is shown in FIG. 2. The opposing electrode 201, the liquid crystals 202, the convex portion 204, and the pixel electrodes 203a and 203b become the structural elements of the simulation model in FIG. 2. The simulation model of FIG. 2 was taken as one unit, and periodically repeated.\nThe simulation parameters are as follows:\ncell gap, d: 4.5 xcexcm, 3.0 xcexcm;\ndistance between pixel electrodes, s: 2 xcexcm, 4 xcexcm;\nconvex portion height, h: 0 xcexcm, 0.2 xcexcm, 0.3 xcexcm, 0.4 xcexcm, 0.5 xcexcm, 0.7 xcexcm, 1.0 xcexcm, and 1.5 xcexcm; and\npixel pitch, p: 18 xcexcm, 43 xcexcm.\nFixed conditions in the simulations are as follows:\nwidth of first region of pixel electrode, o: 1.0 xcexcm;\nelectric potential of the pixel electrode 203a: +5 V;\nelectric potential of the pixel electrode 203b: xe2x88x925 V; and\nelectric potential of the opposing electrode 201: 0 V.\nIn order to generalize the relationship between the distance s between pixel electrodes, the cell gap d, and the height h of the convex portion on which the edge portion of the pixel electrode is formed, the simulation was performed with cell gaps of 4.5 xcexcm and 3.0 xcexcm. ZLI4792 was used for the liquid crystals for both cell gaps, 4.5 xcexcm and 3.0 xcexcm, and the orientation was found by computation.\nThe pre-tilt angle of the liquid crystals was set to 6.0xc2x0, and the chiral pitch was set to left handed at 70 xcexcm. The rubbing direction 205 and 206 are shown in FIG. 2. The twist angle is 90xc2x0. The orientation of the liquid crystal is TN method.\nFurther, in order to increase the number of evaluations and grasp the tendencies, the simulations were performed as stated above with two pixel pitches.\nLiquid crystal orientation simulation software from Syntech Corp. entitled LCD Master 2SBENCH was used, and simulations of the liquid crystal orientation ware run. 2SBENCH shows the liquid crystal orientation by a two-dimensional planer surface composed of the cell cap direction and the substrate surface direction.\nThe simulation results are as shown below. FIGS. 3 to 8 are partial blowup diagrams of the simulation results.\nFIGS. 3 to 8 are simulations in which the height of the convex portion was changed at conditions of the gap d set to 4.5 xcexcm, the inter-pixel spacing distance s set to 2.0 xcexcm, and the pixel pitch p set to 18 xcexcm. The transmittivity in each coordinate, calculated from equipotential lines, liquid crystal directors, and index of refraction anisotropy, is shown. The pixel electrode 203a has coordinates from 1 to 17 xcexcm, and the pixel electrode 203b has coordinates from 19 to 35 xcexcm, and FIGS. 3 to 8 show blowups of the points at which there is light leakage and disclination in the vicinity of the convex portion. The liquid crystal is a positive type, and therefore lines of electric force can be considered to have nearly the same direction as the director of the liquid crystals.\nThe orientation of the liquid crystals on a pixel electrode having a xe2x88x925V electric potential is explained below.\nA horizontal electric field is formed up through a region entering the inside of the pixel electrode when there is no convex portion, as shown in FIG. 3. Further, disclinations develop in regions in which the direction of the horizontal electric field and the direction of the liquid crystal pre-tile angle are opposite.\nWhen the height of the convex portion is 0.3 xcexcm, the vertical direction electric field becomes stronger due to the first region of the pixel electrode formed in the upper edge portion of the convex portion, and therefore the position of the disclination moves toward the outside of the pixel electrode, as shown in FIG. 4, compared with the case of no convex portion in FIG. 3.\nIf the height of the convex portion becomes taller at 0.7 xcexcm, as shown in FIG. 5, then there are many liquid crystals which become oriented perpendicular to the substrate along the lines of electric force due to the effect of the horizontal direction electric field becoming stronger, and due to the effect of the electric force lines, due to the second region of the pixel electrode formed in the side portion of the convex portion and the opposing electrode, possessing components nearly perpendicular with respect to the substrate surface. The disclination in the vicinity of the convex portion becomes less.\nWhen the height of the convex portion becomes higher at 1.0 xcexcm, the width of the disclination in the vicinity of the convex portion drops by just 0.2 xcexcm compared with the case in which the convex portion height is 0.7 xcexcm, as shown in FIG. 6.\nAs shown in FIG. 7, the lines of electric force formed by the second region of the pixel electrode formed in side portion of the convex portion, and by the opposing electrode, possess an angle order of 60xc2x0, with respect to the substrate surface, in the vicinity of the convex portion, which becomes at 1.5 xcexcm. The liquid crystals orient along the electric force lines, and light leakage is formed near the convex portion.\nNote that the angle of the electric force lines with respect to the substrate surface is estimated from the distribution of equipotential lines.\nIn FIG. 8, the height of the convex portion increases to 3.0 xcexcm, and the vertical direction electric field becomes stronger. The liquid crystals on the third region of the pixel electrode of the edge portion of the convex portion therefore orient nearly perpendicular to the surface of the substrate. However, the electric field formed between the second region of the pixel electrode formed in the side portion of the convex portion, and the substrate surface, possesses an angle on the order of 30xc2x0, with respect to the substrate surface, and the liquid crystals orient along the lines of electric force. Broad light leakage is consequently formed in the vicinity of the convex portion.\nThe orientation of the liquid crystals of FIG. 4 corresponds to the schematic diagram of FIG. 1A. The orientation of the liquid crystals of FIG. 5 and FIG. 6 corresponds to the schematic diagram of FIG. 1B, and the orientation of the liquid crystals of FIG. 7 and FIG. 8 corresponds to the schematic diagram of FIG. 1C. Namely, it is confirmed that the light leakage of the liquid crystals becomes larger if the height of the convex portion exceeds an upper limit.\nIn order to improve the display quality, a systematic simulation was performed. The distance between one edge portion and another edge portion of the disclination and the light leakage influencing the aperture ratio was paid attention to.\nData was also taken regarding the light leakage and the disclination width of a pixel electrode having a xe2x88x925 V electric potential. This is because the light leakage and the disclination of a pixel electrode having a xe2x88x925 V electric potential greatly influences the display quality due to high light strength.\nFIG. 9 and FIG. 36 show the simulation results. FIG. 9 is a figure in which the height h of the convex portion, the light leakage and the disclination width x are graphed with respect to the cell gap d in the simulation model of FIG. 2. The term light leakage and disclination width x denotes the width of a region having high brightness due to disclination and light leakage formed on both sides of the convex portion.\nFIG. 36 is a figure in which the height h of the convex portion, and the light leakage and the disclination width y are graphed with respect to the cell gap d in the simulation model of FIG. 2. The term light leakage and disclination width y denotes the width of a region having high brightness due to disclination and light leakage formed in one side of the convex portion, namely in an electrode side having a xe2x88x925 V electric potential.\nSimulations were performed with the pixel pitch p set to 18 xcexcm and 43 xcexcm when the cell gap was 4.5 xcexcm. Further, the distance s between adjacent pixel electrodes was set to 2.0 xcexcm or to 4.0 xcexcm.\nWhen the cell gap was 3.0 xcexcm, the distance s between adjacent pixel electrodes was set to 2.0 xcexcm or 4.0 xcexcm. The pixel pitch p was set to 18 xcexcm.\nThe relationship between the convex portion height, and the width of the light leakage and the disclination shows similar tendencies in both FIG. 9 and FIG. 36.\nFirst, the relationship between the convex portion height and the light leakage and disclination width nearly did not change in accordance to pixel pitch. This is because formation of light leakage and disclination is a phenomenon caused by the horizontal direction electric field and the vertical direction electric field of the edge portions of the pixel electrodes.\nFurther, light leakage and disclination become relatively less with a smaller distance between adjacent pixel electrodes.\nThe regions of orientation irregularities of the liquid crystals, typically light leakage and disclination, are reduced along with increasing height of the pixel portion, and are not due to the pixel pitch p and the cell gap d, in both FIG. 9 and FIG. 36. If the height of the convex portion becomes too tall, there are conversely more regions of orientation irregularities. An optimal convex portion height is determined by the distance between the cell gap and the pixel electrode.\nConsidering the inflection point of the graph, it is preferable that the height of the convex portion in which the effect of a reduction in the liquid crystal regions having orientation irregularities appears significantly, be equal to or greater than 4.4% of the cell gap, and equal to or less than 22.5% of the cell gap when the call gap is 4.5 xcexcm.\nWhen the cell gap is 3.0 xcexcm as well, and the distance s between the pixel electrodes is equal to or less than 2.0 xcexcm, a good effect of reducing, compared to a case in which there is no convex portion, the regions having orientation irregularities can be obtained by setting the height of the convex portion to be equal to or greater than 4.4% of the cell gap, and equal to or less than 22.5% of the cell gap.\nIf the height of the convex portion is less than 4.4% with respect to the cell gap, then the width of the light leakage and the disclination does not change much, even if the height of the convex portion is increased. If the height of the convex portion exceeds 22.5% to the cell gap, then the light leakage and the disclination width increases.\nFurther, orientation irregularities of the liquid crystals easily develop due to rubbing non-uniformity if the convex portion is high, and therefore reducing the height of the convex portion, thereby reducing the light leakage and the disclination width, is preferable to insure display quality. Consequently, for a case in which the cell gap is 4.5 xcexcm, the height of the convex portion may be suppressed to be greater than or equal to 4.4% of the cell gap, and less than or equal to 15.6% of the cell gap. Nearly equal light leakage and disclination reduction effects are obtained when the height of the convex portion is with the range of greater than or equal to 4.4%, and less than or equal to 22.5%, of the cell gap.\nFurther, when the cell cap is 3.0 xcexcm and the distance s between the pixel electrodes is equal to or less than 2.0 xcexcm, a good effect of reducing the light leakage and the disclination is obtained when the height of the convex portion is greater than or equal to 4.4% of the cell gap, and less than or equal to 15.6% of the cell gap, the same as when the height of the convex portion was set equal to or greater than 4.4%, and equal to or less than 22.5%, of the cell gap.\nConversely, there was more light leakage when the height of the convex portion was 22.5% when the distance s between the pixel electrodes was set to 4.0 xcexcm. Inclusive of when the distance between the pixel electrodes is 4.0 xcexcm, it is therefore preferable that the height of the convex portion be equal to or greater than 4.4%, and equal to or less than 15.6%, or the cell gap.\nIn other words, under conditions of the distance between the pixel electrodes being equal to or less than 4.0 xcexcm, the height of the convex portion may be set equal to or greater than 4.4%, and equal to or less than 22.5%, of the cell gap when the cell gap is equal to or greater than 3.0 xcexcm, and equal to or less than 4.5 xcexcm. The height of the convex portion is preferably set greater than or equal to 4.4%, and less than or equal to 15.6%, of the cell gap.\nThe smaller the cell gap, the smaller the height of the convex portion needed for reducing the light leakage and the disclination width. For cases of the call gap being from 3.0 xcexcm to 4.5 xcexcm, good liquid crystal orientation can be obtained with a convex portion height of 15.6% or less. For cases of the cell gap being less than 3.0 xcexcm, it is therefore considered that the necessary convex portion height will be 15.6% or less to sufficiently reduce the light leakage and disclination width.\nThe height of the convex portion may be set equal to or less than 15.6% of the cell gap when the cell gap is less than or equal to 3.0 xcexcm. Of course, considering the inflection point of the graph, it can be expected that a good effect will also be obtained if the height of the convex portion is set equal to or less than 6.7% of the cell gap.\nWhen the cell gap is 3.0 xcexcm, the light leakage and the disclination width are monotonically reduced as the convex portion becomes higher, provided that the height of the convex portion is equal to or less than 6.7% with respect to the cell gap. If the cell gap is made smaller at less than or equal to 3.0 xcexcm, it can be considered that the region in which the light leakage and the disclination width are monotonically reduced with increasing convex portion height is in a range in which the height of the convex portion does not exceed 6.7% with respect to the cell gap.\nAn upper limit to the height of the convex portion, or both upper and lower limits, can thus be determined. There is a concern that rubbing irregularities may occur if the fiber tips of a rubbing cloth are disordered, and therefore the determination of an upper limit for the height of the convex portion is necessary in order to manufacture a liquid crystal panel. Further, the optimal value of the height of the convex portion with respect to the cell gap tends to become smaller as the cell gap becomes smaller.\nAn optimal value for the height of the convex portion thus determined can be used not only for the TN method, but can also be widely used as means for hiding liquid crystal disclination in a normally white mode orientation method.\nThe optimal value of the convex portion height is one in which the lines of electric force formed by a horizontal direction electric field and a vertical direction electric field of an active matrix liquid crystal display device are suitably regulated, and as shown in FIG. 1B, one in which regions with generated electric force lines possessing components perpendicular to the surface of the substrate are increased in the edge portions of the pixel electrode.\nTherefore, although the simulations were performed for a transmitting type liquid crystal display device, it can be considered possible to apply the present invention to a reflective type liquid crystal display device as well. This is because a voltage may also be applied to the pixel electrodes with a reflective type liquid crystal display device, and when orienting the liquid crystals by a vertical direction electric field, unnecessary electric fields directed diagonally with respect to the substrate surface are reduced, and light leakage and disclination of the edge portions of the pixel electrodes can be reduced.\nFurther, the simulation was made using the TN method, but the liquid crystal orientation method is not limited to the TN method. This is because, in an active matrix liquid crystal display device, unnecessary electric fields directed diagonally with respect to the substrate surface can be reduced by optimizing the convex portion height when orienting the liquid crystals by using a vertical direction electric field. For example, it is considered that it is possible to apply the present invention to methods such as an OCB (optically controlled birefringence) method, an STN method, and an EC method using homogeneously oriented cells.\nFurther, provided that orientation faults are not induced in the liquid crystals by the convex portion, it is thought that it is also possible to apply the present invention to an orientation method using sumectic liquid crystals. For example, it is possible to apply the present invention to a liquid crystal display device using ferroelectric liquid crystals or anti-ferroelectric liquid crystals. Further, by adding high molecular weight molecules with liquid crystal properties to these liquid crystals, it is thought that it is also possible to apply the present invention to liquid crystal display devices using materials hardened by irradiation of light (for example, ultraviolet light).\nThe simulations were performed with an angle 90xc2x0 (hereafter referred to as a convex portion taper angle) formed between the surface contacting the second region of the pixel electrode formed in the side surface of the convex portion, and the third region of the pixel electrode formed in the level surface. However, it is also possible to apply the present invention if the convex portion taper angle is less than 90xc2x0. As shown in the cross sectional diagram of FIG. 35 for the lines of electric force when the convex portion has a taper, the electric force lines are generated in a direction perpendicular to a conductor for cases in which the taper angle xcex8 of the convex portion 204 is less than 90xc2x0, and therefore the bend of electric force lines 218 formed between the opposing electrode 201 and second regions 219 of the pixel electrodes 203a and 203b becomes relaxed when the convex portion has a taper. Liquid crystals 220 orient very well perpendicular with respect to the substrate surface. A very large effect in reducing light leakage and disclination is therefore obtained, compared to when the taper angle is 90xc2x0, by using the relationship shown in the present invention between a convex portion height 221 and the cell gap when the convex portion has a taper.\nWidth of the first region of the pixel electrode on the upper portion of the convex portion\nNext, changes in the orientation of liquid crystals when the width of the first region of the pixel electrode formed overlapping with the upper portion of the convex portion is investigated.\nA simulation model is shown in FIG. 2. The opposing electrode 201, the liquid crystals 202, the convex portion 204, and the pixel electrodes 203a and 203b become the structural elements in FIG. 2.\nThe simulation parameters are as follows:\ncell gap, d: 4.5 xcexcm;\ndistance between pixel electrodes, s: 2 xcexcm, 4 xcexcm;\nconvex portion height, h: 0 xcexcm, 0.5 xcexcm; and\nwidth of first region of pixel electrode, o: xe2x88x921 xcexcm, xe2x88x920.5 xcexcm, 0 xcexcm, 0.5 xcexcm, 1.0 xcexcm.\nThe symbol for the width o of the first region of the pixel electrode, such as xe2x88x921.0 xcexcm, indicates that the pixel electrode is not formed on the convex portion, and that the edge portion of the pixel electrode is located at 1.0 xcexcm from the convex portion.\nFixed conditions in the simulations are as follows:\nelectric potential of the pixel electrode 203a: +5 V;\nelectric potential of the pixel electrode 203b: xe2x88x925 V;\nelectric potential of the opposing electrode 201: 0 V; and\npixel pitch, p: 18 xcexcm.\nSimulation results are shown in the cross sectional diagrams of FIG. 11 to FIG. 15. The distance s between pixel electrodes is 2.0 xcexcm.\nIn FIG. 11, there is no convex portion. The convex portion and the pixel electrode do not mutually overlap in FIG. 12, and the edge of the pixel electrode is located 0.5 xcexcm from the edge of the convex portion. In other words, the second region of the pixel electrode and the first region of the pixel electrode do not exist in FIGS. 11 and 12. Light leakage from the edges of the pixel and the disclination width x shows no change at all at this point in FIG. 11 and FIG. 12.\nThe pixel electrode is formed in the side portion of the convex portion in FIG. 13. Namely, there is a second region of the pixel electrode. Compared to FIGS. 11 and 12, the position of the disclination on a pixel electrode having a xe2x88x925 V electric potential, moves 0.4 xcexcm to the pixel edge. The vertical direction electric field is made stronger, and the horizontal direction electric field becomes a little weaker due to the second region of the pixel electrode.\nThe first region of the pixel electrode formed in the edge portion on the convex portion, and the second region of the present invention formed in the side portion of the convex portion exist in FIG. 14. The width of the first region of the pixel electrode is 0.5 xcexcm. The vertical direction electric field becomes stronger, and the disclination on the pixel electrode having a xe2x88x925 V electric potential moves to the edge of the pixel electrode due to the first region of the pixel electrode.\nThe width of the first region of the pixel electrode is set to 1.0 xcexcm in FIG. 15, compared to 0.5 xcexcm in FIG. 14. The vertical direction electric field is made additionally stronger with respect to the horizontal direction electric field due to the 1.0 xcexcm width, and the disclination on the pixel electrode having a xe2x88x925 V electric potential moves to the edge of the pixel electrode.\nIt can thus be understood that there is a disclination reduction effect due to the existence of the first region of the pixel electrode and the second region of the pixel electrode.\nNext, data is added for when the distance s between the pixel electrodes is 4.0 xcexcm, and data is systematically taken. FIGS. 16A and 16B show the simulation results. FIG. 16A is a graph of the width o of the first region of the pixel electrode, and the light leakage and the disclination width x, versus the cell gap d, in the simulation model of FIG. 2. The term light leakage and disclination width x here indicates the width of a region having high brightness causes by disclination and light leakage formed in both sides of the convex portion.\nFIG. 16B is a graph of the width o and the light leakage and the disclination width y, versus the cell gap d, in the simulation model of FIG. 2. The term light leakage and disclination width y here indicates the width of a region having high brightness causes by disclination and light leakage formed in one side of the convex portion, namely in the electrode side having an electric potential of xe2x88x925 V.\nFrom FIGS. 16A and 16B, it can be understood that there is an effect in which disclination and light leakage are reduced if the width o of the first region of the pixel electrode is equal to or greater than 0.5 xcexcm, preferably equal to or greater than 1.0 xcexcm, without dependence on the distance s between the pixel electrodes.\nThe light leakage and the disclination width when the width of the first region of the pixel electrode is 0 xcexcm in FIGS. 16A and 16B shows a light leakage and disclination width of a state in which the pixel electrode is only formed in the side surface of the convex portion. Compared to cases in which the width of the first region of the pixel electrode is equal to or greater than 0.5 xcexcm, or equal to or greater than 1.0 xcexcm, the effect of reducing the light leakage and the disclination width is reduced. However, compared to the case in which the width of the first region of the pixel electrode is xe2x88x920.5 xcexcm, in which the pixel electrode does not contact the convex portion at all, the light leakage and the disclination width do drop.\nAn actual experiment in which the width over which the convex portion and the pixel electrode overlap was carried out. FIG. 33A is an upper surface diagram of a substrate having a convex portion, and FIGS. 33B and 33C are cross sectional diagrams of the substrate having the convex portion.\nPixel electrodes 301a denoted by slanted line portions in the upper surface diagram of FIG. 33A all have the same electric potential. Further, the pixel electrodes 301b denoted by vertical line portions all have the same electric potential. This is in order to connect adjacent pixel electrodes by a transparent conducting film 300 having a width of 3 xcexcm. Assuming line inversion driving, an electric potential of +5 V is imparted to the pixel electrodes 301a. In addition, an electric potential of xe2x88x925 V is imparted to the pixel electrodes 301b. A rubbing direction 302 for the substrate having a convex structure is shown in the figures. A rubbing direction for a substrate opposing the substrate having the convex structure is perpendicular to the rubbing direction 302.\nA cross section of the upper surface diagram of FIG. 33A cut along the dashed line G-Gxe2x80x2 is shown in FIG. 33B. A cross section of the upper surface diagram of FIG. 33A cut along the dashed line H-Hxe2x80x2 is shown in FIG. 33C. Identical reference symbols are used for the same portions as in FIG. 33A. The edge portions of the pixel electrode 301a and 301b formed on the substrate 303 contact a convex portion 304. The distance between the adjacent pixel electrodes 301a and 301b was set constant at 2.0 xcexcm, and the liquid crystal orientation was confirmed by changing the width of the pixel electrode overlapping on the convex portion, namely by changing the width 305 of the first region of the pixel electrode. The width of the first region of the pixel electrode was set to xe2x88x921.0 xcexcm, 0 xcexcm, 0.5 xcexcm, and 1.0 xcexcm. The cell gap was 4.5 xcexcm, the height of the convex portion was 0.5 xcexcm, and the pixel pitch was 18 xcexcm.\nPhotographs of the liquid crystal orientation when the pixel electrode structure of FIGS. 33A to 33C was used are shown in FIGS. 34A to 34D. The adjacent pixel electrodes in the horizontal direction of the page have the same electric potential. The rubbing direction was perpendicular to the page. An effect of reducing the disclination width was found in the experiments as well if the width o of the first region of the pixel electrodes was equal to or greater than 0.5 xcexcm, preferably equal to or greater than 1.0 xcexcm. It was found that the disclination and the light leakage width extending in the horizontal direction of the page decreases along with increases in the width o of the first region or the pixel electrodes."}
{"text": "The present invention relates to a multiple pocket, expandable envelope, a blank for forming a multiple pocket, expandable envelope, and a method of forming a multiple pocket, expandable envelope from a blank. More particularly, the invention relates to the use of a planar, unitary blank for forming a multiple pocket, expandable envelope which may be folded simply and inexpensively in a one pass, right angle gluing and folding operation.\nIn constructing envelopes or folders, it is often necessary to provide a plurality of individual pockets or compartments for separating the contents of the envelope. Additionally, the envelope must be capable of being shipped and stored in a flat, collapsed configuration to use shipping and storage space efficiently, and then expanded to receive varying amounts of contents, e.g., sheets of paper.\nConventional multiple pocket envelopes or folders are formed from a plurality of components which are automatically and manually assembled. The forming of such envelopes from a plurality of components is difficult and expensive, involving both manual steps and complex machinery.\nTo form multiple pocket, expandable envelopes simply, inexpensively and efficiently, such envelopes must be formed from a planar, unitary blank which may be folded and glued in one pass, right angle gluing and folding operation. Such operation eliminates manual handling and permits the envelope to be manufactured on relatively simple, conventional folding and gluing apparatus."}
{"text": "The present invention relates to manufacture of knife blades, and in particular relates to accurate mechanically-performed grinding of sharp edges on a series of similar knife blades.\nIn many knives the blade is tapered between its back and its edge by being beveled on both sides, either with a flat surface or hollow ground. It is desired in sharpening such a blade to provide a narrow edge bevel area on each side of the knife blade, ground at a somewhat steeper angle to provide the actual sharpened edge. It is desired for that narrow portion to conform closely to the profile of the blade, and for the two sides of the blade to be ground symmetrically. The edge bevel areas should be parallel with the entire knife profiles from the plunge to the tip of the blade to be cosmetically appealing.\nMass produced knife blades have slight dimensional variations resulting from the grinders used to produce the flat or hollow ground beveled side surfaces of such blades. In order to produce a desired sharp edge with symmetrical edge bevels on opposite sides of a knife whose blade has been beveled or hollow ground on both sides, it has therefore been necessary in the past for a person to hold each blade and to move the blade manually along a suitable abrasive wheel to sharpen knife blades by hand during the normal process of manufacture. That is, it has not been practical to use machines to automatically grind a final bevel surface to define a sharp edge with the necessary angle to produce a desired appearance, bevel width, and symmetry.\nSome knife blade blanks intended to have serrations in an edge have been beveled flat or hollow-ground on each side to establish a tapered thickness of the blade, narrowing from the back of the blade toward the edge. The tang of such a tapered blade blank has then been clamped into a fixture, and a serrated portion of an edge has been ground into one beveled face using an abrasive wheel dressed to the shape required to provide the desired serrated edge profile. This method of holding a knife blade has not been found useful in sharpening entire knife blades during manufacture, however, because a knife blade cannot be located precisely enough by clamping the tang of the blade to form symmetrical edge bevels reliably on both sides of a blade without having to confirm or readjust the position of the blade between grinding steps. As a result, automatic sharpening of tapered knife blades has not been practical, and such blades have previously been costly because of the amount of skilled labor needed to finish the edge of each blade.\nIt is therefore desired to be able to grind both sides of a knife blade automatically to produce a symmetrically ground sharpened edge without having to readjust or check the blade's position in a fixture holding the blade during the process of grinding the edge."}
{"text": "This invention relates generally to financial services and products, and more particularly to financial systems that enable an investor to create bond laddering strategies, and other strategies that typically involve individual fixed-income securities, through a series of open-ended funds (such as exchange traded funds or mutual funds) that hold fixed-income securities.\nBond laddering is an investment technique in which an investor purchases multiple bonds (or other types of fixed-income securities) having different maturity dates. As an example, the total investment is divided equally among several maturity dates that are spread at regular intervals over a period of time (e.g., every year for a period of ten years). In this way, the bonds associated with each particular maturity date are a “rung” of the total investment “ladder.” Bond laddering has a number of benefits for investors. For example, an investor may avoid capital losses (or gains) due to fluctuations in interest rates by simply holding onto the bonds until the maturity dates arrive, at which time the investor can purchase a new rung for the bond ladder. By having a series of investments that mature over regular intervals, a bond ladder provides more stability around yields and periodic liquidity for the investor.\nA bond laddering strategy may also minimize risks to the investor. For example, in an environment where interest rates are increasing, a bond ladder prevents the investor from being stuck with a low interest rate for a long time, thereby allowing the investor to benefit from the returns associated with increasing rates. On the other hand, when interest rates are decreasing, the division of the investment across several rungs of the bond ladder means that only one portion of the total bond investment (i.e., one rung of the ladder) will mature at any given time. Accordingly, if interest rates are relatively low but the investor desires to repurchase more bonds when a rung matures, only a fraction of the total portfolio will consistent of the bonds at the lower-yielding rate. As one can appreciate, by spreading out the maturity dates of the fixed-income securities over time, the ladder structure smoothes the risks associated with investing in fixed-income securities in an environment where interest rates can fluctuate.\nInvestors commonly invest in different types of funds to gain exposure to various types of securities, including bonds and other fixed-income securities. A very popular type of fund is an exchange-traded fund, or ETF. Shares of an ETF are securities that represent a legal right of ownership over an underlying portfolio of securities or other assets held by the issuing fund. The assets held by an ETF may include individual stocks, bonds, cash, commodities, derivatives, or any tradable asset, including contracts based on the value of any of the foregoing. Shares of an ETF are designed to be listed on a securities exchange and traded over the exchange just like other securities. ETFs thus allow an investor to own a set or “basket” of assets by simply purchasing shares in the individual ETF. Many existing ETFs hold a mix of assets that aim to replicate or otherwise match the characteristics of a particular published index. These ETFs allow investors to have exposure to the index by purchasing shares of the single ETF. Because of their low cost and tax advantages, ETFs have grown in popularity in recent years.\nA typical ETF resembles an index mutual fund in that an ETF generally holds a basket of securities designed to replicate the returns of a securities index and is required to permit daily redemptions at the current value of its holdings (also known as “Net Asset Value”). But unlike mutual funds, ETF shares trade on an exchange throughout the trading day, and most investors buy and sell shares on the exchange rather than direct purchases and redemptions from the fund itself (as is the case with mutual funds). Unlike mutual funds, most transactions in ETF shares are conducted in the secondary market (i.e., on an exchange) and do not involve the movement of assets in or out of the fund. In the case of transactions in creation units that do involve the movement of assets into or out of the fund, the transactions are routinely effected by giving the redeeming shareholder their pro rata share of the fund's holdings, which does not impose trading costs or adverse tax consequences on the remaining shareholders.\nAlthough there are many ETFs that hold bonds, existing ETF products are perpetual in nature and therefore do not enable investors to implement strategies that rely on a maturity structure such as bond laddering. Other types of funds that may hold bonds (including open-ended funds, such as mutual funds) similarly do not enable an effective bond laddering strategy. Aside from the lack of a liquidation date, this limitation is, at least in part, due to the open-ended nature of these funds. A bond laddering strategy typically relies on the ability to buy and hold a fixed-income security until its maturity date, thus allowing an investor to lock in a specific yield for each rung of the bond ladder. But with traditional open-ended funds like ETFs, there is no liquidation date and no limitation on investors' ability to enter or leave the fund, which may require the fund to buy or sell the underlying bonds when investors buy or sell shares of the fund. The lack of a liquidation date coupled with the buying and selling of the underlying bonds while interest rates are changing may result in a net change of the yield of the underlying funds and result in uncertainty for investors around the value of their investment at liquidation. These attributes have prevented traditional open-ended funds from creating a constant yield experience for each investor in the fund. As such, existing open-ended funds are unacceptable for strategies which rely upon a liquidation date and expected yield such as bond laddering.\nAccordingly, it would be desirable to provide a fund management information system 110 that solves these problems by offering and managing open-ended funds (such as ETFs) that have properties allowing the funds to be used by investors to create and pursue other applications typically associated with individual fixed income securities such as bond laddering strategies."}
{"text": "In recent years, there are many database systems that allow storage and querying of extensible Markup Language data (“XML data”). Though there are many evolving standards for querying XML, all of them include some variation of XPath. However, database systems are usually not optimized to handle XPath queries, and the query performance of the database systems leaves much to be desired. In specific cases where an XML schema definition may be available, the structure and datatypes used in XML instance documents may be known. However, in cases where an XML schema definition is not available, and the documents to be searched do not conform to any schema, there are no efficient techniques for querying using XPath.\nSome database systems may use ad-hoc mechanisms to satisfy Xpath queries that are run against documents where the schema of the documents is not known. For example, a database system may satisfy an XPath query by performing a full scan of all documents. While a full scan of all documents can be used to satisfy all XPath queries, the implementation would be very slow due to the lack of indexes.\nAnother way to satisfy XPath queries involves the use of text keywords. Specifically, many database systems support text indexes, and these could be used to satisfy certain XPaths. However, this technique can only satisfy a small subset of XPath queries. Thus, there have been no effective indexing techniques in existing database systems that can be used to handle a wide variety of XPath queries.\nThe approaches described in this section are approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section."}
{"text": "1. Field of the Invention\nThe present invention is related to a lithography system.\n2. Background Art\nIn lithography systems, pattern generators are used to pattern substrates during an exposure process. Example lithography systems include, but are not limited to, reflective or transmissive maskless, immersion, and mask-based system. Example substrates include, but are not limited to, semiconductor wafers, flat panel display substrates, flexible substrates, and the like. Light interacting with an illumination spot on the pattern generator becomes patterned. The patterned light is projected using a projection optical system onto one or more target areas of the substrate during the exposure process to form features on a photosensitive material (e.g., a photoresist) on the substrate.\nIn some applications of the lithography system, it is desired to have double telecentric illumination, which means that a chief ray of an illumination beam is telecentric in both object space and image space. Object space is proximate an object plane, typically defined by the pattern generator. Image space is proximate an image plane, typically defined by the substrate. Conventional systems having double telecentric illumination generally use a transmissive optic to direct the illumination from an illuminator towards the pattern generator and from the pattern generator towards a projection optical system. However, transmissive optics tend to absorb illumination, and thus can have low efficiency. For example, a typical transmissive optic can sometimes transmit as little as about 20–25% of the light between the illuminator and the projection system. Also, transmissive optics can cause polarization, aberration, distortion, birefringence, and other errors in the illumination, which sometimes requires complex optics to correct.\nTherefore, what is needed is a system and method that can be used to provide a double telecentric system, which eliminates or substantially reduces absorption of illumination and/or that eliminates or substantially reduces polarization change, birefringence, aberration, distortion, and/or other errors being introduced in the illumination beam."}
{"text": "Heretofore some kinds of 4-aryl-4-halopyrazolidine-3,5-dione derivatives have been known. For example, An. Quim. vol. 71, 396 to 399 (1975) discloses that compounds A, B and C to be described later, i.e. 4-aryl-4-halopyrazolidine-3,5-dione derivatives, can be prepared by Diels-Alder reaction between 4-chloro-4-phenylpyrazoline-3,5-dione and diene. Further An. Quim. vol. 75, 931 to 935 (1979) describes the under-mentioned compounds D and E as an intermediate for preparing 13,14-azasteroid, the compounds being prepared by Diels-Alder reaction between 4-chloro-4-phenylpyrazoline-3,5-dione and diene. In addition, Comment., Pontif. Acad. Sci. vol.2, 8 (1975) sets forth that compounds A, B and C to be described later can be prepared by chlorination of the corresponding 4-phenylpyrazolidine-3,5-dione derivative.\nHowever, the compounds of the invention represented by the formula (1) are novel compounds which are not included in conventional compounds.\nOn the other hand, 4-aryl-4-nitropyrazolidine-3,5-dione derivatives are entirely novel compounds. ##STR2##\nWith insecticides and miticides in use for many years, pest insects have acquired resistance in recent years and become difficult to control with conventional insecticides. Accordingly it is expected to develop novel insecticides. Also in the herbicide field, the development of novel herbicides are expected.\nAn object of the present invention is to provide 4-aryl-4-substituted pyrazolidine-3,5-dione derivatives which have not only extremely high insecticidal and miticidal activities but a high herbicidal activity."}
{"text": "While many food containers have lids that snap or “lock” onto the container to retain the food inside the container (e.g., Tupperware® and VERSAtainer™), these containers do not secure the food from unauthorized access or tampering.\nThe issues of tampering and access have been addressed in the prior art with respect to bottles for medications and the like through the use of child-proof caps and single-use tamper-indicating rings, tabs, foils and the like, as well as locking medicine cabinets.\nWhile medications have been determined to be worthy of such means for protection from tampering and access, the prior art has not considered meal dishes to be worthy of protection from tampering and unauthorized access."}
{"text": "The invention pertains to methods of forming dielectric materials, and methods of forming capacitors. The invention also pertains to capacitor constructions.\nIt is frequently desired to form dielectric materials during semiconductor device fabrication. For instance, capacitor constructions comprise dielectric material separating a pair of capacitor electrodes. Suitable dielectric materials for capacitor constructions include silicon dioxide and silicon nitride, with an exemplary dielectric material comprising a stack of silicon nitride between a pair of silicon dioxide layers.\nAn advantage of utilizing silicon nitride in capacitor constructions is that it has a higher dielectric constant than silicon dioxide. However, a difficulty in utilizing silicon nitride can occur in attempting to get a uniform coating of silicon nitride over a capacitor electrode. For instance, a capacitor electrode can comprise conductively-doped rugged silicon (for example, conductively-doped hemispherical grain polysilicon). Such rugged silicon has a rough surface texture, and is utilized because the rough surface texture enables more conductive surface area to be provided over a particular footprint than would be provided with a smooth-surfaced structure. A difficulty can occur in attempting to form silicon nitride over the roughened surface structure of rugged silicon. Specifically, silicon nitride is typically provided by chemical-vapor deposition (such as, for example, low pressure chemical vapor deposition utilizing silane and ammonia as precursors), and the nitride deposits non-conformally on the roughed surface of the rugged silicon. Accordingly, if the nitride is provided as a thin layer (less than 100 xc3x85 thick), there can be pinholes extending into the nitride, and even extending through the nitride to expose portions of the underlying rugged silicon surface.\nAmong the methods which have been developed to compensate for the pinhole problems are methods in which silicon dioxide is formed over the layer of silicon nitride to either fill the pinholes or at least cover the pinholes with a dielectric material. The silicon dioxide can be formed by either chemical vapor deposition, or by oxidation of the is silicon nitride surface.\nAnother method for compensating for pinhole problems is to form silicon dioxide over the rugged polysilicon prior to formation of the silicon nitride. Accordingly, a dielectric material will be beneath the silicon nitride, and any pinholes extending through the silicon nitride can be prevented from exposing the underlying conductive substrate by the intervening layer of silicon dioxide.\nIn typical prior art processing, both of the above-discussed silicon dioxide methodologies are utilized. In other words, a layer of silicon dioxide is formed before forming the layer of silicon nitride, and a second layer of silicon dioxide is formed after forming the layer of silicon nitride.\nIt would be desirable to develop methods wherein some or all of the above-discussed difficulties associated with formation of silicon nitride could be eliminated, and particularly it would be desirable to develop methods wherein one or both of the above-discussed layers of silicon dioxide could be eliminated from capacitor constructions.\nIn one aspect, the invention encompasses a method of forming a dielectric material. A nitrogen-comprising layer is formed on at least some of the surface of a rugged polysilicon substrate to form a first portion of a dielectric material. After the nitrogen-comprising layer is formed, at least some of the substrate is subjected to dry oxidation with one or both of NO and N2O to form a second portion of the dielectric material.\nIn another aspect, the invention encompasses a method of forming a capacitor. A layer of rugged silicon is formed over a substrate, and a nitrogen-comprising layer is formed on the layer of rugged silicon. Some of the rugged silicon is exposed through the nitrogen-comprising layer. After the nitrogen-comprising layer is formed, at least some of the exposed rugged silicon is subjected to dry oxidation conditions with one or both of NO and N2O. Subsequently, a conductive material layer is formed over the nitrogen-comprising layer.\nIn yet another aspect, the invention encompasses a capacitor structure. The structure includes a first capacitor electrode comprising a rugged polysilicon layer, a nitrogen-comprising layer on the rugged polysilicon layer, and a second capacitor electrode. The nitrogen-comprising layer is between the first and second capacitor electrodes."}
{"text": "This invention relates to the use of ether amide derivatives as low calorie fat mimetics. These compounds have a (C.sub.2 to C.sub.12) backbone to which are attached one or two aliphatic groups in ether linkage and one to four aliphatic groups in amide linkage.\nDietary fat is the most concentrated source of energy of all the nutrients, supplying 9 kcal/gram, about double that contributed by either carbohydrate or protein. The amount of fat in the American diet has increased in the last 60 years by about 25% (Mead, J., et al. Lipids, Plenum, N.Y., 1986, page 459), so that fats now provide approximately 40% of the daily caloric intake. Moreover, technological advances in the food industry, including efficient and safe hydrogenation procedures, have changed the kind of fat in foods.\nBecause fats are high in calories and because certain fats appear to pose a health risk when consumed in large quantities over time, a number of national advisory committees on nutrition have made recommendations differing in detail, but the common theme is a reduction in the total amount of fat in the diet (Gottenbos, J. J., chapter 8 in Beare-Rogers, J., ed., Dietary Fat Requirements in Health and Development, A.O.C.S. 1988, page 109). Yet fat contributes to the palatability and flavor of food, since most food flavors are fat-soluble, and to the satiety value, since fatty foods remain in the stomach for longer periods of time than do foods containing protein and carbohydrate. Furthermore, fat is a carrier of the fat-soluble vitamins, A, D, E, and K, and the essential fatty acids, which have been shown to be important in growth and in the maintenance of many body functions. Hence, major research efforts have focused on ways to produce food substances that provide the same functional and organoleptic properties as fats, but not the calories.\nA number of fat replacements have heretofore been suggested (recently reviewed by Hamm, D. J., J. Food Sci. 49: 419-428 (1984), Haumann, B. J., J. Amer. Oil Chem. Soc. 63: 278-288 (1986) and LaBarge, R. G., Food Tech. 42: 84-90 (1988)). Hamm divides replacement fats into two broad categories: structurally re-engineered triglycerides modified to retain their conventional functional properties in foods, while removing their susceptibility toward hydrolysis or subsequent absorption during digestion, and materials developed from chemistry unrelated to triglycerides.\nExamples of the former class of triglyceride analogues include compounds having the glycerol moiety replaced with alternate polyols (e.g., pentaerythritol in U.S. Pat. No. 2,962,419 to Minich, or sugars, suggested by Hess, K., and Messmer, E., Ber. 54B: 499-523 (1921), and patented years later by Mattson and Volpenhein, U.S. Pat. No. 3,600,186, and Meyer, et al., U.S. Pat. No. 4,840,815); compounds having the fatty acids replaced with alternate acids (e.g., branched esters as described in U.S. Pat. No. 3,579,548 to Whyte); compounds having insertions between the glycerol and the fatty acid (e.g., ethoxy or propoxy groups in U.S. Pat. No. 4,861,613 to White and Pollard); compounds having reversed esters (e.g., malonates in U.S. Pat. No. 4,582,927 to Fulcher and trialkoxytricarballylates in U.S. Pat. No. 4,508,746 to Hamm); and compounds having the ester bonds replaced by ether bonds (Can. Pat. No. 1,106,681 to Trost).\nExamples of Hamm's second category of fat replacements chemically unrelated to triglycerides are mineral oil (suggested as early as 1894 in U.S. Pat. No. 519,980 to Winter); polyglucose and polymaltose (U.S. Pat. No. 3,876,794 to Rennhard); jojoba wax (W. Ger. Pat. No. 3,529,564 to Anika); polyoxyalkylene esters (U.S. Pat. No. 4,849,242 to Kershner); polyvinyl alcohol esters (U.S. Pat. No. 4,915,974 to D'Amelia and Jacklin); and polysiloxane (Eur. Pat. Ap. No. 205,273 to Frye).\nNondigestible or nonabsorbable edible fat replacements have proved disappointing when tested in feeding trials, where gastrointestinal side effects occurred, in some cases so extreme that frank anal leakage was observed. Nondigestible fats appear to act as a laxative and are expelled from the body, eliciting foreign body reactions like those early documented for mineral oil (Stryker, W .A., Arch. Path. 31: 670-692 (1941), more recently summarized in Goodman and Gilman's Pharmacological Basis of Therapeutics, 7th ed., Macmillan Pub. Co., N.Y. 1985, pp. 1002-1003). Similarly, a series of experimental fats, e.g., glyceride esters of dibasic acids, synthesized by U.S.D.A. in the 1960's exhibited undesirable gastrointestinal side effects when the compounds were fed to rats (Booth, A. N., and Gros, A. T., J. Amer. Oil Chem. Soc. 40: 551-553 (1963)); in several of the balance studies, the diarrhea was so extreme that digestibility coefficients could not be calculated (ibid., Table I, p. 552).\nPolyglycerol and polyglycerol esters, suggested as fat replacements by Babayan and Lehman (U.S. Pat. No. 3,637,774), have been suggested for use as fecal softening agents as well (U.S. Pat. No. 3,495,010 to Fossel). A number of remedies have been recommended to combat the anal leakage observed when sucrose polyesters are ingested (e.g., employing cocoa butters, U.S. Pat. No. 4,005,195 to Jandacek, incorporating saturated fatty groups, Eur. Pat. Ap. No. 233,856 to Bernhardt, or mixing residues, U.S. Pat. No. 4,797,300 to Jandacek, et al.), and dietary fiber preparations have been incorporated into polysaccharide and/or polyol-containing foodstuffs to help inhibit the diarrheal effect (U.S. Pat. No. 4,304,768 to Staub et al.). Partially digestible fat replacements have also been suggested (U.S. Pat. No. 4,830,787 to Klemann and Finley; U.S. Pat. No. 4,849,242, cited above; and U.S. Pat. No. 4,927,659 to Klemann, et al.)."}
{"text": "It is highly desirable for tires to exhibit good traction characteristics on both dry and wet surfaces. However, it has traditionally been very difficult to improve the traction characteristics of a tire without compromising its rolling resistance and tread wear. Low rolling resistance is important because good fuel economy is virtually always an important consideration. Good tread wear is also an important consideration because it is generally the most important factor which determines the life of the tire.\nThe traction, tread wear and rolling resistance of a tire is dependent to a large extent on the dynamic viscoelastic properties of the elastomers utilized in making the tire tread. In order to reduce the rolling resistance of a tire, rubbers having a high rebound have traditionally been utilized in making the tire's tread. On the other hand, in order to increase the wet skid resistance of a tire, rubbers which undergo a large energy loss have generally been utilized in the tire's tread. In order to balance these two viscoelastically inconsistent properties, mixtures of various types of synthetic and natural rubber are normally utilized in tire treads. For instance, various mixtures of styrene-butadiene rubber and polybutadiene rubber are commonly used as a rubber material for automobile tire treads. However, such blends are not totally satisfactory for all purposes.\nRubbers having intermediate glass transition temperatures (-70.degree. C. to -40.degree. C.) compromise rolling resistance and treadwear without significantly increasing traction characteristics. For this reason, blends of rubbers having low glass transition temperatures and rubbers having high glass transition temperatures are frequently utilized to attain improved traction characteristics without significantly compromising rolling resistance or treadwear. However, such blends of rubbers having low glass transition temperatures and rubbers having high glass transition temperatures exhibit poor processability. This major disadvantage associated with such blends has greatly hampered their utilization in making tire tread compounds.\nTin-coupled polymers are known to provide desirable properties, such as improved treadwear and reduced rolling resistance, when used in tire tread rubbers. Such tin-coupled rubbery polymers are typically made by coupling the rubbery polymer with a tin coupling agent at or near the end of the polymerization used in synthesizing the rubbery polymer. In the coupling process, live polymer chain ends react with the tin coupling agent thereby coupling the polymer. For instance, up to four live chain ends can react with tin tetrahalides, such as tin tetrachloride, thereby coupling the polymer chains together."}
{"text": "A portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent document or the patent disclosure as it appears in the Patent and Trademark Office patent file or records, but otherwise reserves all copyright rights whatsoever. The following notice applies to the software and data as described below and in the drawings referred to herein: Copyright 2005, Cyberscan Technology, Inc., All Rights Reserved."}
{"text": "The present invention relates to a process of thermal treatment in a tissue, and especially relates to a process of thermal treatment in a tissue comprising embedding a heat-sensitive magnetic material (hereinafter referred to as \"implant\") in a tumor, and externally applying a high frequency magnetic field thereto for raising the temperature of the implant in order to promote the medical treatment.\nA process of thermal treatment in a tissue is one in which high frequency induction heating is utilized for medically treating a tumor in a living body in which a plurality of needle-like implant elements are embedded, based on the therapeutic effects which destroy only cancer cells to minimize the unfavorable influence on normal cells when the temperature of tissues of an organism is elevated to 42.degree. C. to 43.degree. C. The process of thermal treatment in a tissue is supposed to be promising not only for a brain tumor but also for other internal organ cancers such as a tongue cancer, a spleen cancer, a bladder cancer, a bone tumor, a mammary cancer and [Tatsuya Kobayashi, Yoshihisa Kida, Masaaki Matsui and Yoshifumi Amamiya, NEUROLOGICAL SURGERY, Vol. 18, No. 3 (issued on Mar. 10, 1990)]. A metallic alloy such as a Ni-Cu alloy, a Ni-Si alloy, a Ni-Pd alloy, an Fe-Pt alloy, an Fe-Pd alloy and an Fe-Pt-Si alloy and an oxide such as .gamma.-hematite, magnetite, cementite, rare earth magnetic materials, strontium ferrite, magnet-plumbite type ferrite and amorphous materials are known as material of the implant element employed in the process of thermal treatment (Japanese patent laid open gazette No. 63-103048).\nHeat generation of the above implant material of the metallic magnetic material is carried out by mainly utilizing eddy current loss, while that of the magnetic material made of an oxide or that of the powders of metallic materials are carried out by mainly utilizing magnetic hysteresis loss of powders. When the eddy current is utilized, the magnetic field is required to be applied in the longitudinal direction of a cylindrical sample for avoiding the influence of a demagnetizing field. For avoiding the directional influence of the magnetic field, the magnetic field may be rotated utilizing a Helmholz type coil which generates a magnetic field by parallel connecting the same sized coils to provide a space in the middle. Such properties as a sufficient amount of heat generation, a relatively low Curie temperature and the safety in an organism are required as implant material. It has been conventionally supposed that an Fe 73% (atomic %, any percentage appearing in the rest of the specification is also atomic %)- Pt 27% alloy having a low Curie temperature of 68.degree. C. is most suitable.\nSince, however, even the temperature of the implant element made of the Fe-Pt alloy itself rises to 68.degree. C., there is a possibility of damaging the normal cells as well as the cancer cells by means of the heat generated. When the implant element made of the Ni-Cu alloy or the Ni-Si alloy having a similar Curie temperature to that of the Fe-Pt alloy is employed, a large number of the implant elements are required to be embedded in an affected part because of the insufficient heat generation. On the other hand, in case of the bulk-like molded material made of the powders utilizing the hysteresis loss, there are disadvantages such that when the high frequency magnetic field is applied, an eddy current flows at a contacted portion to prevent the sufficient penetration of the magnetic flux into the implant element so as to suppress the generation of heat by means of the hysteresis loss that results in the uneven heat transmission to the affected part. Although the amount of the heat generation is desirably changed according to the individual difference of a patient and the progress of the affected part, it is in fact impossible.\nFurther, the compatibility of the implant element to an organism including the dispersion of the component elements into the organism as a result of corrosion and ionization is problematic even if such a metal as Fe and Pt is a metal of considerably low toxicity because the implant element is embedded in the organism for a long period of time. The molded member made of powders such as ceramics becomes brittle after the repetition of the thermal treatments in a tissue so that the uncovered powders are in danger of straying in the living body."}
{"text": "FIELD OF INVENTION\nThis invention relates generally to slatwall-type display shelving and, more specifically, to interlocking slatwall-type panels and method therefor."}
{"text": "It is known that in the deep-drawing of sheets, particularly relatively thin sheets, the frictional situation between the holding down device and the pull part is important, since this influences the continuity of flow of the workpiece into the die. In order to obtain workpieces free from wrinkles and cracks, in addition to the friction other parameters have to be adapted such as the material, drawing part shape, drawing slot, hold-down pressure and the like. For example, wrinkle formation at the drawing part cannot be prevented by unlimited increase of the hold-down pressure, since increased frictional forces result in crack formation in the drawing part. Also the use of lubricants does not result in the desired success, if the usual smooth or directed surface structures are employed. The film of the lubricant in this case is squeezed away from the hold-down device such that the friction at this place of highest load rapidly reaches an impermissible level and results in a tearing of the drawing part. It is also known that an improvement of the frictional situation at the drawing slot and at the hold-down device results if in place of the usual rolled-bright, matte, brushed, ground or with foil-covered flat surface a roughened surface is employed, which allows at the contact surface the production of a high lubricant pressure. Thus it is known to employ work rollers with corrugated valleys, which are generated by grinding, to produce a roughened surface (Austrian patent disclosure No. AT-PS 325,236), or trough shaped impressions which are generated by sand blasting (Austrian patent disclosure No. AT-PS 347,387), for aluminum and aluminum alloys for rolling, since the impressions keep the lubricant in place during the rolling. A sheet produced in this way does not show advantageous properties in the subsequent deep-drawing process, since the valleys are pulled by the rolling operation in longitudinal direction with a large longitudinal elongation and are thereby directed, so that the lubricant can escape in this direction during the drawing process."}
{"text": "1. Field of the Invention\nThe present invention relates to processes for the preparation of optically active amines or salts thereof and processes for the preparation of optically active imines. The invention also relates to novel optically active amines and salts thereof and to novel optically active imines.\n2. Description of the Related Art\nCertain optically active amines are known in the art as intermediates for medical drugs and agricultural chemicals. For example, an optically active 7-amino-5-azaspiro[2.4]heptane is a compound useful as an intermediate for synthetic antimicrobial agents that have high antibacterial activities.\nSuch an optically active 7-amino-5-azaspiro[2.4]heptane has been prepared by certain processes typically described in Japanese Unexamined Patent Application Publications No. 11-240867 and No. 11-240868. In these processes, an oxo group at the 7-position of a precursor 5-benzyl-4,7-dioxo-5-azaspiro[2.4]heptane is oximated with hydroxylamine, and an oxo group at the 4-position and a hydroxyimino group at the 7-position are reduced to yield 7-amino-5-benzyl-5-azaspiro[2.4]heptane, the 7-amino-5-benzyl-5-azaspiro[2.4]heptane is then condensed with an optically active carboxylic acid and is converted into a 5-benzyloxycarbonyl derivative, the resulting diastereomers are separated and are deprotected to thereby yield an optically active 7-amino-5-azaspiro[2.4]heptane.\nThese publications also describe a process for the preparation of an optically active 7-amino-5-azaspiro[2.4]heptane derivative, in which an optically active 1-phenylethyl group is introduced, instead of the benzyl group, at the 5-position of 5-benzyl-4,7-dioxo-5-azaspiro[2.4]heptane, the oxo group at the 7-position of the resulting compound is converted into an amino group, and the resulting diastereomers are separated to thereby yield the target compound.\nIn all of these processes, however, diastereomers are once prepared synthetically and are then subjected to resolution. The processes are thereby low in yields and are not efficient for the preparation of an optically active 7-amino-5-azaspiro[2.4]heptane."}
{"text": "Electrostatic discharge (ESD) is the sudden and momentary electric current that flows between two electric elements at different electrical potentials. ESD events may cause damage to electronic equipment, especially in solid-state electronics, such as integrated circuits. As fabrication of integrated circuits is moving into sub-micron processes, the circuits are more vulnerable to stress due to ESD. That is, as fabrication, e.g., CMOS, processes downscale, ESD robustness degrades due to the metallization layers being thinner, hence having higher sheet resistance.\nRecently developed techniques allow fabricating radio frequency integrated circuits (RFICs) that operate in millimeter wave frequencies, typically using deep sub-micron CMOS technology. In order to enable proper and safe operation of ICs, on-chip ESD protection should be included. In a design of an RFIC, a critical point that requires ESD protection is at the RF front-end between the low noise amplifier (LNA) and antenna(s). The LNA, which is commonly the most sensitive element in the RF front end, is typically exposed to ESD sparks (pulses) induced through the antenna pins.\nConventional ESD protection techniques applicable for ICs cannot be efficiently implemented in RFICs and in particular, RFICs for millimeter wave frequencies. The conventional techniques introduce significant parasitic capacitance and resistance, thus degrading RF front-end performance.\nIn the related art, several on-chip ESD protection techniques for RFIC can be found. For example, the article “ESD Protection for a 5.5 GHz LNA in 90 nm RF CMOS Implementation Concepts, Constraints and Solutions” to Natarajan, et al., proposes to place inductors at input and output pins of the LNA to divert the ESD current away from the LNA core to the power lines. The size of the inductors can be relatively small, thus suitable for RFIC designs where area is limited. However, the fabricated inductors are characterized by high series resistance due to the thin metal layer of the semiconductor substrate. This reduces the effectiveness of the ESD protection. Further, the inductors introduce high heat dissipation in a long-lasting ESD pulse. This can lead, for example, to open interconnect.\nAnother ESD solution designed for protection of RFICs is based on a shorted-stub transmission line. As illustrated in FIG. 1A, the transmission line 110 is connected between an input pin 101 of a LNA 120 and the ground. The shorted-stub transmission line 110 discharges ESD pulses, while at the same time operates as a matching circuit for a normal operation of the LNA 120. This ESD protection solution suffers from the deficiencies noted above, e.g., high resistance and high heat dissipation. In addition, the length of the transmission line 110 is typically λ/4, where λ is the wavelength of a signal amplified by the LNA 120. For example, for 60 GHz signals the length of the transmission line 110 is not less than 600 micron (when fabricating a transmission line on a multilayer semiconductor substrate the length of each line is shorter than λ/4 due to characteristics of the substrate).\nIn order to improve the protection efficiency and reduce the heat dissipation, a ladder-shaped ESD protection circuit is proposed in “Ladder-shaped network for ESD protection of millimeter-wave CMOS ICs” to Park, et al. As illustrated in FIG. 1B, the ESD protection circuit consists of three shorted stubs 130 and three series stubs 140. The length of each shorted stub 130 is 300 micron. The ladder-shaped structure operates as a high-pass filter, i.e., filtering ESD pulses characterized by a low frequency relative to the operating frequency of the LNA.\nThe major disadvantage of the shorted-stub based ESD protection solutions, illustrated in FIGS. 1A and 1B, is that they are considerably large in size for RFIC designs. Therefore, such structures are unsuitable for use in RFICs, and more particularly elements that should be integrated in devices for 60 GHz band applications.\nTherefore, it would be advantageous to provide a solution for designing a compact size circuit for providing robustness in ESD protection for RFICs operation in millimeter wave frequencies."}
{"text": "This invention relates to osteoimplants and their use in the repair of bone defects and injuries. More particularly, the invention relates to a load-bearing composite osteoimplant which can assume any of a wide variety of configurations, methods for their manufacture and the use of the osteoimplants for the repair of hard tissue.\nShaped or cut bone segments have been used extensively to solve various medical problems in human and animal orthopedic surgical practice, and their application has also extended to the field of cosmetic and reconstructive surgery, dental reconstructive surgery, and other medical fields involving surgery of hard tissues. The use of autograft bone (where the patient provides the source), allograft bone (where another individual of the same species provides the source) or xenograft bone (where another individual of a different species provides the source) is well known in both human and veterinary medicine. In particular, transplanted bone is known to provide support, promote healing, fill bony cavities, separate bony elements (such as vertebral bodies), promote fusion (where bones are induced to grow together into a single, solid mass), or stabilize the sites of fractures. More recently, processed bone has been developed into shapes for use in new surgical applications, or as new materials for implants that were historically made of non-biologically derived materials.\nBone grafting applications are differentiated by the requirements of the skeletal site. Certain applications require a xe2x80x9cstructural graftxe2x80x9d in which one role of the graft is to provide mechanical or structural support to the site. Such grafts contain a substantial portion of mineralized bone tissue to provide the strength needed for load-bearing. The graft may also have beneficial biological properties, such as incorporation into the skeleton, osteoinduction, osteoconduction, or angiogenesis.\nStructural grafts are conventionally made by processing, and then cutting or otherwise shaping bones collected for transplant purposes. The range of bone grafts that might be thus prepared is limited by the size and shape limitations of the bone tissue from which the bone graft originated. Certain clinically desirable shapes and sizes of grafts may thus be unattainable by the cutting and shaping processes, due to the dimensional limitations of the bone. For some shapes they may also be available only in limited amounts, due to the large variations inherent in the human or animal donor source populations.\nMany structural allografts are never fully incorporated by remodeling and replacement with host tissue due, in part, to the difficulty with which the host\"\"s blood supply may penetrate cortical bone, and partly to the poor osteoinductivity of nondemineralized bone. To the extent that the implant is incorporated and replaced by living host bone tissue, the body can then recognize and repair damage, thus eliminating failure by fatigue. In applications where the mechanical load-bearing requirements of the graft are challenging, lack of replacement by host bone tissue may compromise the graft by subjecting it to repeated loading and cumulative unrepaired damage (mechanical fatigue) within the implant material. Thus, it is highly desirable that the graft have the capacity to support load initially, and be capable of gradually transferring this load to the host bone tissue as it remodels the implant.\nIt is an object of the present invention to provide an osteoimplant possessing sufficient strength in a body fluid environment to enable the osteoimplant to bear loads.\nIt is a further object of the present invention to provide a load-bearing osteoimplant which contains pores or cavities which permit the osteoimplant to be revascularized and incorporated by the host.\nIt is yet a further object of the present invention to provide a load-bearing osteoimplant which is osteogenic and thereby promotes new host bone tissue formation within and around the osteoimplant.\nIt is yet another object of the invention to provide a load-bearing osteoimplant which supports load initially and is capable of gradually transferring this load to the host bone tissue as it remodels the osteoimplant.\nIt is yet another object of the invention to provide a load-bearing osteoimplant containing a reinforcing component.\nIt is yet an even further object of the present invention to provide methods for the manufacture of osteoimplants of any size and/or configuration ranging from the relatively simple to the relatively complex.\nIt is yet an even further object of the present invention to provide methods for the manufacture of bone-containing osteoimplants which are not limited by constraints imposed by the shape and/or size of the bone tissue from which the osteoimplants are manufactured.\nIt is still another object of the invention to provide tissue-engineered bone manufactured from one or more synthetic materials and tissue obtained from various sources such as transgenic animals, plants and microorganisms.\nIt is yet another object of the invention to provide an integral implant insertion instrument and implant possessing an implant portion in accordance with the invention which, following its implantation in the body, is separated from the insertion portion of the instrument.\nThese and further objects of the invention are obtained by a load-bearing osteoimplant which comprises a shaped, coherent aggregate of bone particles.\nThe load-bearing osteoimplant of this invention is fabricated by the method which comprises providing an aggregate of bone particles, optionally in combination with one or more additional components, and shaping the mass into a coherent unit of predetermined size and shape employing at least one or more processes such as extruding, molding, -solvent/gel casting, machining, computer aided design and computer aided manufacturing (CAD/CAM), and the like.\nThe bone particles utilized in the fabrication of the osteoimplant of this invention are selected from the group consisting of nondemineralized bone particles, demineralized bone particles and combinations thereof. The bone particles are remodeled and replaced by new host bone as incorporation of the osteoimplant progresses in vivo. As described more fully hereinbelow, the bone particles can be fully demineralized by removing substantially all of their inorganic mineral content, they can be partially demineralized by removing a significant amount, but less than all, of their inorganic mineral content or they can be superficially demineralized by confining the removal of inorganic mineral to just the surface of the bone particles.\nThe term xe2x80x9costeoimplantxe2x80x9d as utilized herein contemplates any device or material for implantation that aids or augments bone or other hard tissue formation or healing for human or animal use. Osteoimplants are often applied at a bone defect or dental repair site, e.g., one resulting from injury, defect brought about during the course of surgery, infection, malignancy or developmental malformation. Therefore, osteoimplants are envisioned as being suitably sized and shaped as required for use in a wide variety of orthopedic, neurosurgical, oral and maxillofacial and dental surgical procedures such as the repair of simple and compound fractures and non-unions, external and internal fixations, joint reconstructions such as arthrodesis, general arthroplasty, deficit filling, discectomy, laminectomy, anterior cervical and thoracic operations, spinal fusions, dental restorations, etc. Therefore, the osteoimplants herein are intended for implantation at a bony site and, in addition to bone particles, can contain one or more other components, e.g., binder (adhesive), filler, biologically active component, reinforcing component or reinforcing structure, coupling agent, as described in U.S. Pat. No. 6,399,693, the contents of which are incorporates by reference herein, -and the like.\nThe term xe2x80x9cshapingxe2x80x9d refers to any of the various methods that can be used, individually or in combination, to provide an osteoimplant of a desired size and configuration. Such methods include, for example, extruding, injection molding, solvent casting, vacuum forming, sintering, melt forming, reaction molding, compression molding, transfer molding, blow molding, rotational molding, thermoforming, machining, CAD/CAM procedures, and the like, and include any post-shaping operations that may be utilized to modify the internal and/or external structure of the osteoimplant and/or modify its properties, e.g., selective removal of a filler component to provide voids, application of a layer of biologically active material to part or all of the surface and/or subsurface region of the osteoimplant, bonding of bone particles through the crosslinking of their mutually contacting exposed collagen, etc.\nThe term xe2x80x9cbiocompatiblexe2x80x9d and expressions of like import shall be understood to mean the absence of stimulation of an unacceptable biological response to an implant as distinguished from the sort of mild, transient inflammation and/or granulation response which can accompany implantation of foreign objects into a living organism and which is associated with the normal healing response. Optional components that are useful can be considered biocompatible if, at the time of implantation, they are present in a sufficiently small concentration such that the above-defined condition is achieved.\nThe term xe2x80x9cparticlexe2x80x9d as applied to the bone component of the osteoimplant includes bone pieces of all shapes, sizes, thickness and configuration such as fibers, threads, narrow strips, thin sheets, chips, shards, powders, etc., that posses regular, irregular or random geometries. It should be understood that some variation in dimension may occur in the production of the bone particles and bone particles demonstrating considerable variability in dimensions and/or size are within the scope of this invention. Bone particles that are useful herein can be homogenous, heterogeneous and can include mixtures of human, xenogenic and/or transgenic material.\nThe term xe2x80x9chumanxe2x80x9d as utilized herein in reference to suitable sources of bone refers to autograft bone which is taken from at least one site in the graftee and implanted in another site of the graftee as well as allograft bone which is human bone taken from a donor other than the graftee.\nThe term xe2x80x9cautograftxe2x80x9d as utilized herein refers to tissue that is obtained from the intended recipient of the implant.\nThe term xe2x80x9callograftxe2x80x9d as utilized herein refers to tissue, which may be processed to remove cells and/or other components, intended for implantation that is taken from a different member of the same species as the intended recipient. Thus, the term xe2x80x9callograftxe2x80x9d includes bone from which substantially all cellular matter has been removed (processed acellular bone) as well as cell-containing bone.\nThe term xe2x80x9cxenogenicxe2x80x9d as utilized herein refers to material intended for implantation obtained from a donor source of a different species than the intended recipient. For example, when the implant is intended for use in an animal such as a horse (equine), xenogenic tissue of, e.g., bovine, porcine, caprine, etc., origin may be suitable.\nThe term xe2x80x9ctransgenicxe2x80x9d as utilized herein refers to tissue intended for implantation that is obtained from an organism that has been genetically modified to contain within its genome certain genetic sequences obtained from the genome of a different species.\nThe term xe2x80x9ccompositexe2x80x9d as utilized herein refers to the mixture of materials and/or components used in preparing the shaped osteoimplant.\nThe term xe2x80x9cwholexe2x80x9d as utilized herein refers to bone that contains its full, or original, mineral content.\nThe term xe2x80x9cdemineralizedxe2x80x9d as utilized herein refers to bone containing less than about 95% of its original mineral content and is intended to cover all bone particles that have had some portion of their original mineral content removed by a demineralization process. Non-demineralized bone particles provide strength to the osteoimplant and allow it to initially support a load. Demineralized bone particles induce new bone formation at the site of the demineralized bone and permit adjustment of the overall mechanical properties of the osteoimplant.\nThe expression xe2x80x9cfully demineralizedxe2x80x9d as utilized herein refers to bone containing less than about 8% of its original mineral context.\nThe term xe2x80x9costeogenicxe2x80x9d as utilized herein shall be understood as referring to the ability of an osteoimplant to enhance or accelerate the growth of new bone tissue by one or more mechanisms such as osteogenesis, osteoconduction and or osteoinduction.\nThe term xe2x80x9costeoinductivexe2x80x9d as utilized herein shall be understood to refer to the ability of a substance to recruit cells from the host that have the potential for forming new bone and repairing bone tissue. Most osteoinductive materials can stimulate the formation of ectopic bone in soft tissue.\nThe term xe2x80x9costeoconductivexe2x80x9d as utilized herein shall be understood to refer to the ability of a non-osteoinductive substance to serve as a suitable template or substrate along which bone may grow.\nThe term xe2x80x9costeoimplantxe2x80x9d herein is utilized in its broadest sense and is not intended to be limited to any particular shape, size, configuration or application.\nThe term xe2x80x9cshapexe2x80x9d as applied to the osteoimplant herein refers to a determined or regular form or configuration in contrast to an indeterminate or vague form or configuration (as in the case of a lump or other solid mass of no special form) and is characteristic of such materials as sheets, plates, disks, cores, pins, screws, tubes, teeth, bones, portions of bones, wedges, cylinders, threaded cylinders, cages, and the like. . This includes forms ranging from regular, geometric shapes to irregular, angled, or non-geometric shapes, and combinations of features having any of these characteristics.\nThe term xe2x80x9cimplantablexe2x80x9d as utilized herein refers to a biocompatible device retaining potential for successful surgical placement within a mammal.\nThe expression xe2x80x9cimplantable devicexe2x80x9d and expressions of like import as utilized herein refers to any object implantable through surgical, injection, or other suitable means whose primary function is achieved either through its physical presence or mechanical properties.\nThe term xe2x80x9cbioresorbablexe2x80x9d as utilized herein refers to those materials of either synthetic or natural origin which, when placed in a living body, are degraded through either enzymatic, hydrolytic or other chemical reactions or cellular processes into by-products which are either integrated into, or expelled from, the body. It is recognized that in the literature, the terms xe2x80x9cresorbablexe2x80x9d, xe2x80x9cabsorbablexe2x80x9d, and xe2x80x9cbioabsorbablexe2x80x9d are frequently used interchangeably.\nThe term xe2x80x9cpolymericxe2x80x9d as utilized herein refers to a material of natural, synthetic or semisynthetic origin that is made of large molecules featuring characteristic repeating units.\nThe expression xe2x80x9calternating copolymersxe2x80x9d as utilized herein refers to copolymers with a regular or alternating repeating unit sequence The expression xe2x80x9cthermoplastic elastomersxe2x80x9d as utilized herein refers to melt-processable copolymers which possess elastomeric mechanical properties as a result of a crystallizable xe2x80x9chardxe2x80x9d segment and an amorphous xe2x80x9csoftxe2x80x9d segment possessing a Tg below its service temperature.\nThe term xe2x80x9cblendsxe2x80x9d as utilized herein refers to polymeric materials that are melt-mixed to achieve compounding between two or more different polymeric compositions that are not covalently bonded to each other. For the purposes of this application, a melt-miscible blend is a polymeric mixture that possesses sufficient miscibility in the melt to be useful in shaping.\nThe term xe2x80x9cincorporationxe2x80x9d utilized herein refers to the biological mechanism whereby host cells gradually remove portions of the osteoimplant and replaces the removed portions with native host bone tissue while maintaining strength. This phenomenon is also referred to in the scientific literature by such expressions as xe2x80x9ccreeping substitution,xe2x80x9d xe2x80x9cwound healing responsexe2x80x9d and xe2x80x9ccellular based remodeling.xe2x80x9d Therefore, the term xe2x80x9cincorporationxe2x80x9d shall be understood herein as embracing what is considered by those skilled in the art to be conveyed by the aforequoted expressions.\nThe expression xe2x80x9cfurther treatmentxe2x80x9d as utilized herein refers to procedures such as lyophilization, cross-linking, re-mineralization, sterilization, etc., performed either before, during or after the step of shaping the bone particle-containing aggregate as well as post-process procedures such as machining, laser etching, welding, assembling of parts, cutting, milling, reactive etching, etc. It further includes treatment(s) applied at the time of surgery such as rehydration, combining with cellular materials, application of growth factors, etc.\nThe expression xe2x80x9cwet compressive strengthxe2x80x9d as utilized herein refers to the compressive strength of the osteoimplant after the osteoimplant has been immersed in physiological saline (water containing 0.9 g NaCl/100 ml water) for a minimum of 12 hours and a maximum of 24 hours. Compressive strength is a well known measurement of mechanical strength and is measured using the procedure described herein."}
{"text": "Crude oil is a complex mixture of organic and inorganic species. Pitch is the residue that is left after refining crude oil under vacuum and high temperatures. It is generally considered waste and is the major component of road tar. Pitch contains thousands of aromatic hydrocarbons with 3 to 8 fused ring systems, which is also known as carbon mesophase. The carbonaceous pitches containing asphaltene are useful in the formation of a variety of carbon artifacts due to their high carbon-to-hydrogen ratio. There is a need for additional methods and resulting compositions that can use pitch to generate useful asphaltene and asphaltene derivatives."}
{"text": "As the need for increased dam storage changes, the search for higher density, faster access memory technologies also increases. One of these, holographic data storage, provides the promise for increased access to higher density data. The techniques for realizing such storage typically utilize some type of storage media, such as photorefractive crystals or photopolymer layers, to store 3-D \"stacks\" of data in the form of pages of data. Typically, coherent light beams from lasers are utilized to perform the addressing, writing and reading of the data from the storage media by directing these beams at a specific region on the surface of the media. Writing is achieved by remembering the interference pattern formed by these beams at this region. Reading is achieved by detecting a reconstructed light beam as it exits the storage medium, the data then being extracted therefrom. Addressing is achieved by the positioning of the laser beams, and this is typically done through the mechanical movement of mirrors or lenses; however, the storage media itself can be moved relative to fixed laser beams.\nOne method for positioning the laser beam onto the storage media is to utilize a fixed optics system that selects a given storage location, such that both the reference beam and the data beam are focused onto the storage location. In order to select another storage location on a two-dimensional holographic media, the storage media can then be moved along its x- and y-axes. Alternatively, the optic system can be more complicated, such that the beams can be located to a particular region without the need to move the media itself. However, apparatus for achieving such movement are quite complex."}
{"text": "Field of the Invention\nThe present invention relates to an image forming apparatus, a drum unit, and a developing unit. More specifically, the present invention relates to an image forming apparatus, a drum unit, and a developing unit to form a color image by toners of plural colors.\nDescription of Related Art\nAn image forming apparatus for forming a color image by a tandem electrophotographic color printer includes an image carrier and a developing roller for each color to form a toner image with the corresponding color. The image carrier is used up by image forming, and hence a drum unit including the thus consumed image carrier needs to be replaced. Further, when the drum units have different configurations for each color, some measures are required to prevent wrong insertion of a drum unit for different color for replacement.\nFor example, Japanese Patent Application Publication No. JP 2010-217658 teaches a technique of providing protrusions and recesses in a drum unit and in an attachment part of an image forming apparatus in which the drum unit is to be attached in order to prevent wrong attachment of the drum unit. Specifically, these protrusions and recesses are provided in different positions for each color, and thus any drum unit other than a correct drum unit for the corresponding color is not allowed to be attached to the attachment part.\nHowever, some drum units are not so largely different in visual appearance among different colors. Therefore, there is a case that a user erroneously inserts a wrong drum unit into an attachment part for the different color.\nFurther, the protrusions and the recesses for preventing the wrong attachment are formed by resin molding, which is an inexpensive method in general, and they do not have much strength. As a result, when a user erroneously inserts the drum unit in the attachment part for wrong color, the protrusions and the recesses for preventing the wrong attachment may be broken.\nThe present invention has been made to solve the above problem of the conventional technique. To be specific, the present invention has objects of providing an image forming apparatus enabling to prevent wrong attachment of a drum unit without causing breakage when a replacement part has the wrong attachment, and further providing a drum unit and a developing unit which can be used as the replacement part of the image forming apparatus."}
{"text": "The present invention relates to a plurality of feeder aids mounted on a feeder drum of a self-propelled combine harvester.\nA first feeder drum is disclosed in EP 0 836 800 B1 that is composed of feeder segments mounted on an axle in rows. The feeder segments are essentially square and are mounted on the axle such that they are offset relative to each other and are non-rotatable. A wide threshing cover is located on every corner of the feeder segment for reinforcement.\nThe disadvantage of this configuration of the feeder drum is its high-cost, laborious production outlay, since every segment must be manufactured individually and mounted on the axle. In addition, the covers, with their wide configuration, have a negative effect on the threshing of the crops, because the wide effective surface area of the covers creates high friction on the crops and, due to this friction, a high portion of damaged grain often results. Finally, when harvesting conditions are unfavorable, the fact that the crops are acted upon by a large surface area sometimes causes an entire ear to be torn off, thereby continually impairing the separation of the grain."}
{"text": "A turbine wastegate is typically a valve that can be controlled to selectively allow at least some exhaust to bypass a turbine. Where an exhaust turbine drives a compressor for boosting inlet pressure to an internal combustion engine (e.g., as in a turbocharger), a wastegate provides a means to control the boost pressure."}
{"text": "The background art contains several examples of methods for constructing fibrous mats. An example is disclosed in U.S. Pat. No. 5,565,049 to Simmons, et al. Simmons discloses a method of making mats from chopped fibrous material. The method uses an elongated endless transport belt that extends over a pair of rollers. At an initial station, a series of spools deliver continuous fiber strands into a cutting station. The cut strands are deposited onto the belt to form a loose fibrous mat structure. The mat structure is then transported to a misting station where nozzles spray a fine mist of water onto the mat structure. An adhesive powder is then applied uniformly over the top surface of the mat structure. The powdered adhesive takes the form of granules that are metered out of an adhesive applicator. At a subsequent heating station an upper belt is provided that applies pressure to the mat. The heating station also includes heating elements positioned above and below the mat structure. The pressure and heat cause the mat structure to flatten out and the adhesive to become plasticized. Next, a cooling station is provided that brings the temperature of the mat down below the initial plasticizing temperature of the adhesive. This causes the adhesive to re-solidify and bind the fibers together.\nAnother method of making fibrous mats is disclosed in U.S. Pat. No. 5,051,122 to Reese et al. Reese discloses a method of making mats of continuous glass fiber strands in which a heated calendaring roll and continuous belt are used to compact the mat. A plurality of strand feeders are traversed across the surface of a moving perforated conveyor with each feeder drawing at least one continuous strand from a supply source and projecting it onto the surface of the conveyor belt to form a loose fibrous mat structure. Powdered resin particles are distributed over the loose mat, which is then passed between a second moving conveyor or belt and at least one heated calendaring roller. One embodiment uses water to wet the mat for easier distribution of the resinous material. As the mat passes to a second conveyor belt it is pinched between the surfaces of the conveyor and a separately driven, rotating, heated calendar roll. As the mat passes around the circumference of the roll it is compacted against it by the tension of the belt. As the belt turns, it separates the compacted mat from the calendar roll and continues its path. Steam from the heating of the mat by the calendar roll is vented away by means of a hood.\nAlthough each of the above references discloses a useful construction method, the references also suffer from certain disadvantages. For instance, each of the references applies a resin or adhesive in a particulate or granular form. This application technique requires multiple steps and does not adequately coat the underlying fibers. As a result there is insufficient bonding between adjacent fibers. The background art likewise fails to illustrate a method for constructing fibrous mats in conjunction with composite building boards. The system and method of the present disclosure are aimed at overcoming these and other shortcomings in the background art."}
{"text": "The present invention relates to a process for treating complex cyanide, especially cyanide associated with spent aluminum potliners.\nAluminum potliners are created in the smelting of aluminum metal and its alloys. They are the liners of the pots into which the molten aluminum is poured. A typical liner lasts about five years. The Pacific Northwest aluminum producers currently annually produce about 50,000-75,000 tons of spent potliners, which are currently classified in the U.S. as hazardous waste. They contain significant concentrations of several soluble pollutants, namely, cyanides, fluorides, and heavy metals (such as lead, beryllium, and cadmium). They may also include polynuclear aromatic hydrocarbons (PNAs).\nU.S. patent application Ser. No. 09/254,000 (which I incorporate by reference) describes a treatment process for spent potliners that is technically feasible at a reasonable cost (with reasonable capital cost). That process produces disposable waste products that comply with all environmental control regulations. It is robust and able to tolerate variations in the feed because spent potliners from different sources differ significantly in their makeup (that is, the process can accept widely differing feedstreams rather than requiring a particular feedstock). The process initially reacts crushed potliner with a calcium chloride/hydrochloric acid (CaCl2/HCl) leach mill solution at elevated temperatures and pressures. Then, it oxidizes the PNAs with a strongly acidic FeCl3 leach mill solution in an oxygen-containing environment at elevated temperature and pressure. Cyanides are destroyed and complex fluorides converted to insoluble florspar by reacting a slurry of finely crushed potliner (xe2x85x9 inch minus) in the calcium chloride/hydrochloric acid leach mill solution at a pulp density of about 30-40%.\nComalco treats potliners using a two-step process involving a calciner to complex the cyanides followed by a hydrometallurgical step to recover the fluorides. Pechiney\"\"s SPLIT process introduces ground potliner to a hot airflow vortex rotating at supersonic speeds to destroy the cyanides followed by reacting the fluorides with an additive to produce a disposable solid waste.\nThe xe2x80x9cLCLL process, xe2x80x9d Alcan\"\"s Low Caustic Leaching and Liming process, digests (leaches) finely ground spent potliner in a dilute caustic solution at around 85xc2x0 C. for about one hour. Then, Alcan adds NaOH in a plugflow reactor at about 180xc2x0 C. and 160 psig to destroy the cyanide in the leach solution in about one additional hour of processing while producing sodium fluoride. Finally, Alcan adds more caustic (generally, lime) to the remaining fluoride liquor for a third hour of treatment in equipment comparable to aluminum smelter wet scrubbers to produce calcium fluoride and a recyclable, caustic leach solution.\nThe Adrien process for treating spent potliner uses five stages to recover calcium fluoride, aluminum fluoride suitable for smelting, and to leave a carbonaceous, solid waste. Adrien leaches the crushed potliner, washes the solid residue with two acid washes using NMF or another fluoroacid, and, finally, uses a water wash.\nThe process of the present invention destroys complex cyanides effectively and efficiently. In air at atmospheric pressure an effective amount of chloride salt solution, especially calcium chloride, magnesium chloride, or a mixture of these salts, and sodium or calcium hypochlorite (i.e., household bleach) react within about one hour with complex cyanide, such as cyanide complexed with iron or copper as is common with aluminum spent potliner, to reduce the cyanide. Optionally (and typically), the process also uses an effective amount of calcium carbonate as a source of carbon dioxide to accelerate the reaction, as illustrated in the formulae below. I believe that the following reactions are involved with my preferred process (using calcium salts for purposes of illustration):\nCaCl(OCl)+CO2+H2O=CaCl(HCO3)+HClOxe2x80x83xe2x80x83(1)\nHClO=HCl+Oxe2x80x83xe2x80x83(2)\n2HCl+CaCO3=CaCl2+CO2+H2Oxe2x80x83xe2x80x83(3)\nThe combination of ingredients and reactants causes the cyanide to oxidize. I have not been able to destroy the cyanides in potliner using chloride salt solution alone or the hypochlorite alone, at least when the cyanide is bound with iron or copper. The chloride salt and hypochlorite together, however, provides an effective environment for destruction of the cyanides. The calcium carbonate forms carbon dioxide that appears to accelerate the destruction of the cyanide."}
{"text": "A digital transport stream typically comprises scrambled audio/video/data content packets and scrambled conditional access messages (ECM, EMM) which have to be processed in order to extract control data (CW); the latter being necessary for the multimedia unit to descramble the content packets relating to the event (e.g a broadcasted program or a set of programs). Conditional access messages include two kinds of information, namely Entitlement Control Messages (ECM) and Entitlement Management Messages (EMM). The ECM is information relating to content packets, whereas the EMM is information dedicated to an individual end user (subscriber) or a group of end users. The ECM contains the access conditions for the current broadcast event together with control word (CW) for unscrambling this event. The Control Word (CW) is a key used for decrypting content packets of a packetized event broadcasted within the data stream. Thus, each ECM is specific to each event (e.g. a TV program). The EMM is a message used typically by the security module to set, reset or change product access entitlements, credit, etc. . . . Thus EMM refers to the rights (access data) of the subscriber for accessing to the content (events).\nThere is two main types of set-top-boxes known by the prior art. The first type comprises set-top-boxes that can generally receive broadcast unencrypted events, and which therefore do not require any access control. As internet access is widespread at user homes, such set-top-boxes are increasingly provided with an IP-connection which allows the end user to use his set-top-box together with his television screen as a terminal for internet facilities.\nThe second type refers to set-top-boxes provided each with a security module. These set-top-boxes are typically used by Pay-TV systems for processing selective access to broadcast services on a per-subscriber basis managed by a conditional access system (CAS). Such a security module typically refers to a smart card (chip card) which can be inserted into a slot of a conditional access reader. This card reader can be connected to the multimedia unit or be lodged directly within this unit. Thus, the security module becomes a necessary device for decrypting the scrambled content packets of the digital transport stream. To this end, the conditional access messages contained in the transport stream are routed to the security module which has the keys to decrypt the ECM in order to extract the control word (CW) which then will be used by a descrambling module (within the multimedia unit) for descrambling the content packets. This process is achieved only if the conditional access rights (provided by the EMM) are valid and checked by the security module.\nThe problem with set-top-boxes of the first type is that they are often unsuitable for processing conditional access messages required for descrambling audio/video content packets within a Pay-TV environment. This is true even if they are provided with an IP-connection. Thus, a user wanting to subscribe a subscription for having access to Pay-TV programs has to change his set-top-box with another one that supports conditional access processing.\nOn the other hand, the set-top-boxes of the second type, provided with a security module, are dependent on this module for descrambling the content packets of the transport stream. It means that the security module of the set-top-box is required for descrambling purposes. If this security module becomes unavailable (for instance because the required smart card has to be changed to a renewal or for any other reasons), this set-top-box becomes at least temporarily unavailable for receiving Pay-TV programs.\nAccordingly, there exists a need for improving processes and multimedia units (e.g. set-top-boxes) in order to optimize the processing of control data for descrambling the content packets of the transport stream, for instance by adapting this processing to different circumstances, depending on the availability and the adequacy of the network and/or depending on the material configuration at the end user."}
{"text": "Bioresin is a term coined to describe a resin or resin formulation derived from a biological source. Thus, many traditional resins such as protein-based soybean, collagen or casein, or carbohydrate-derivatives from cellulose or starch, natural rubber based adhesives and natural phenolic adhesives such as tannin or lignin may all be classed as bioresins. They are renewable polymers and, with low environmental impact, represent an alternative to existing petroleum-driven systems.\nInterest in bio-resin systems stems largely from increasing regulation and public concern for a pollution-free environment, and the need for sustainable alternatives to products based upon a finite petrochemical resource. Commercial thermoset resin production and use is subject to such regulation largely due to the monomer components that form the basis of the formulation. Overwhelmingly, these resins are based on melamine, phenol, urea, formaldehyde, styrene or isocyanate starting materials.\nEmissions of formaldehyde, phenol and isocyanate are generally regulated; they are governed in England and Wales under the Pollution Prevention and Control Regulations (2000), SI19, for installations involved in the manufacture of particleboard (PB), oriented strandboard (OSB), medium density fibreboard (MDF) and wood fibreboard. Emissions to air are limited to 5 mg/m3. In particular, the guidance notes that ‘operators should use resins which minimise emissions of formaldehyde wherever possible. The choice of resins used should be continually reviewed to ensure minimum emissions occur. Problems may exist concerning both high formaldehyde levels in the workplace and with the slow release of formaldehyde from the panel products themselves, particularly at the beginning of the product's lifetime.\nBioresin products may offer an alternative, renewable source of thermosetting resins that will begin to address the depletion of finite resources and have better emissions profiles, particularly when formulated without formaldehyde. Though formaldehyde-free materials such as iisocyanate and tannin resins have become available, about 85% of MDF produced today uses formaldehyde resin, corresponding to a European market of approximately 2 million tonnes per annum. In those resins, whilst a reduction in formaldehyde content has been achieved in recent years, it has been at the cost of longer processing times, and decreases in internal bond and bending strengths and an increase in swelling and water absorption of panel products made therefrom. There is, thus, a clear place in the market now for new more environmentally friendly resins that are competitive in price, performance and adaptable to existing composite manufacturing processes."}
{"text": "In a cell array of a conventional semiconductor memory device, a memory cell may include a switching element and a data storage element, and a cell may have an 8F2 configuration (F being a feature size based on a design rule). Many efforts have been made to provide increased integration of devices. More particularly, an area occupied by a memory cell has been reduced. For example, semiconductor memory devices are changing from 8F2 configurations to 6F2 configurations. Active regions of a 6F2-type cell array recently developed may be diagonal.\nPortions of a conventional cell array having diagonal active regions is illustrated in FIG. 1 and FIG. 2. The conventional cell array has a plurality of diagonal active regions. Adjacent active regions are arranged face to face in a major axis direction. To define diagonal shapes on a reticle, a to-be-exposed portion is divided into segments S and S′ of a predetermined width which are exposed discontinuously. Therefore, more time may be required to fabricate a reticle than a square configuration having vertical and horizontal axis directions. Further, a segment size may be determined according to pattern-to-pattern distances B and B′ and a pattern inclination so that correcting patterns may be difficult. If patterns are disposed face to face in a major direction, a bridge may occur at patterns transcribed to a substrate if a distance between adjacent patterns is a feature size A. In the case where patterns are diagonally disposed face to face, as illustrated in FIG. 2, a pattern-to-pattern distance A′ may become narrower than a feature size increasing a possibility that a bridge may occur.\nTo form a normal pattern during an exposure process, an optical requirement for a light source of a direction corresponding to an inclined axis of a pattern may be provided. Thus, an optical system may adopt a unique aperture to select and irradiate light in the direction corresponding to a pattern axis."}
{"text": "This disclosure relates generally to database management systems and, more particularly, relates to transforming a user-defined table function (UDTF) to a derived table in a database management system (DBMS). Databases are used to store information for numerous types of applications. Examples include various industrial, commercial, technical, scientific, and educational applications. Database management systems (DBMSs) are a typical mechanism for accessing data stored in a database. DBMSs are typically configured to separate the process of storing data from accessing, manipulating, or using data stored in a database. DBMSs often require tremendous resources to handle the heavy workloads placed on such systems. As such, it may be useful to transform a UDTF to a derived table in a DBMS."}
{"text": "The present invention relates to a method and apparatus for the take-up of continuous filament bundles of synthetic fibers at a high speed. More particularly, it relates to a method for the take-up of continuous filament bundles of synthetic fibers, such as yarns and tows, or the like at a high speed, comprising introducing the continuous filament bundles carried on the air stream into a bent tube provided downstream of the air jetting device, which is rotating in substantially reverse direction to that of the continuous filament bundles discharged from the outlet of the bent tube at a peripheral speed of the outlet thereof being 0.5 to 2.0 times the take-up speed of the continuous filament bundles, and receiving the continuous filament bundles into a receiving device, and an apparatus therefor.\nHitherto, the take-up of a number of yarns (e.g. tow) has been carried out by using a pair of gear rollers or a pair of belts, between which the tow is put and taken up. With the speed up in the treatment of the tow for improving the productivity and for rationalizing the procedure, it is required to take up the tow at a higher speed. For instance, when the tow is taken up at a high speed of 2,000 m/min. or more, the conventional methods using gear rollers or belts have a defect that the tow is wound around the final take-up rollers due to the accompanying air stream. For eliminating the defect, it has been proposed to provide an air jetting device in the high speed take-up method. However, when the tow discharged from the air jetting device is directly received into the tow receiving device, the received tow is disturbed by the air stream jetted out together with the tow, and thereby, when the tow is drawn out from the receiving device in the subsequent step, the tow becomes entangled and can not stably be drawn out. There has also been proposed a method comprising striking the tow discharged from the air jetting device onto a lattice or net drum, by which the tow is separated from the air stream, and then the tow is taken up. However, according to this method, when the tow is struck, it is somewhat crimped and entangled, and therefore, when the tow is drawn out from the receiving device in the subsequent step, the tow can not stably be drawn out because it has become snarled.\nUnder the circumstances, the present inventors have intensively attempted to eliminate these defects in the conventional methods for taking up the continuous filament bundles. As the results, it has been found that the object can be accomplished by using a bent tube which is rotating in a substantially reverse direction to that of the continuous filament bundles discharged from the outlet of the bent tube at a fixed speed.\nAn object of the present invention is to provide an improved take-up method of the continuous filament bundles of synthetic fibers at a high speed.\nAnother object of the invention is to provide an apparatus useful for the take-up of the continuous filament bundles of synthetic fibers at a high speed.\nThese and other objects will be apparent from the following description.\nThe take-up method of the continuous filament bundles of synthetic fibers (hereinafter, referred to merely as \"yarns\") of the present invention comprises introducing the yarns discharged from the air jetting device and carried on the air stream into a bent tube provided downstream from the air jetting device, which is rotating in substantially reverse direction to that of the yarns discharged from the outlet of the bent tube at a peripheral speed of the outlet thereof being 0.5 to 2.0 times the take-up speed of the yarns, and receiving the yarns into the receiving device provided downstream from the bent tube.\nThe apparatus useful for the take-up of yarns at a high speed according to the present invention comprises an air jetting device provided downstream from the final take-up rollers, a bent tube provided downstream from the air jetting device, and a receiving device provided downstream from the bent tube.\nAccording to the present invention, the winding of the yarns onto the final take-up rollers can be prevented by using the air jetting device, and further, by introducing the yarns discharged from the air jetting device into the bent tube together with the jetting air stream and moving the tip of the bent tube circularly in substantially a reverse direction to that of the stream of the yarns and the air, at a speed approximately equal to the take-up speed of the yarns, the yarns and the jetting air are discharged from the outlet of the bent tube at an absolute speed of approximately zero. Therefore the yarns are stably, without any disturbance, disposed into the receiving device."}
{"text": "The invention relates to a method and an apparatus for the cleaning and compacting of metal and/or ceramic material, such as lathe cuttings and/or grinding grits.\nIn the machining or grinding of workpieces large amounts of cuttings or grits are produced which are wet with emulsions or with oils from the machining process. The oil or the emulsion is necessary to the working process. These liquids serve to cool the workpiece and to carry away the cuttings. The cuttings are separated from the liquids in the filters of the machines in question. The liquid itself remains in the machines, while the grits and/or cuttings have to be removed from the filters.\nDue to the great surface area of the cuttings, large amounts of lubricant and coolant come out with them. To reduce this loss of liquid the cuttings are centrifuged or pressed. The liquid is returned into the circuit, while the remaining residue must be disposed of. Disposing of these cuttings and/or grits presents considerable difficulty for environmental reasons.\nSince the cuttings and/or grits are contaminated with mineral oils, esters, glycerol, emulsifiers, additives containing phosphorus or heavy metals etc., these useful substances should not be recycled directly into the melting process. Another problem is the low bulk density of the cuttings. For proper disposal the cuttings and/or grits would have to be decontaminated and thoroughly compacted.\nIn the known cleaning processes the degreasing is performed by washing, pressing or centrifugation. In thermal processes the heat is applied by indirect inductance, resistance heating or even flame heat sources."}
{"text": "A three-dimensional quadrupole ion trap includes three electrodes which define a chamber. Two of the three electrodes are virtually identical and, while having hyperboloidal geometry, resemble small inverted saucers. The electrodes which resemble inverted saucers are called end-cap electrodes and are typically distinguishable by a number of holes in the center of each electrode. For example, one end-cap electrode may have a single small central aperture through which ions can be gated periodically, and the other end-cap electrode may have several small centrally arranged apertures through which ions can be ejected from the chamber of the ion trap so as to interact with a detector. (Note that ion traps which utilize external ion sources typically have a single perforation in each end-cap electrode.) The third electrode also has hyperboloidal geometry and is called the ring electrode. The ring electrode is positioned symmetrically between the two end-cap electrodes, and all three cooperate to define the aforementioned ion trap chamber.\nThe geometries of the electrodes are defined so as to produce a quadrupole field which, in turn, will produce an ion trapping potential for the confinement of ions in an area within the chamber of the ion trap defined by the ion trapping potential. For example, an ion trapping potential can be created from a field generated when an oscillating potential is applied to the ring electrode and the two end-cap electrodes are grounded.\nBecause a quadrupole ion trap can generate an ion trapping potential for the confinement of ions, it can function as an ion storage device in which gaseous ions can be confined for a period of time in the presence of a buffer gas, such as 1 mTorr of helium gas. For example, as a storage device, the ion trap can act as an “electric field test-tube” for the confinement of gaseous ions, either positively or negatively charged, or both, in the absence of solvent.\nOne use of the confinement of gaseous ions in such a “test-tube” permits the study of gas-phase ion chemistry. In addition, the ion trap can also function as a mass spectrometer in that the mass-to-charge ratios of the confined ions can be measured. For example, as each ion species is ejected from the chamber of the ion trap in a mass selected fashion, the ejected ions impinge upon an external detector thereby creating a series of ion signals dispersed in time which constitutes a mass spectrum. Ejection of ions from the chamber of the ion trap can be accomplished by ramping, in a linear fashion, the amplitude of a radio frequency (r.f.) potential applied to the ring electrode; each ion species is ejected from the chamber (and thus the area defined by the ion trapping potential) at a specific r.f. amplitude and, because the initial amplitude and ramping rate are known, the mass-to-charge can be determined for each ion species upon ejection. This method for measuring mass-to-charge ratios of confined ions is known as the “mass-selective axial instability mode”.\nOne area of interest in which the above described ion traps are utilized is the study of large polyatomic molecules such as peptides, proteins, oligonucleotides, carbohydrates, and synthetic polymers. These polyatomic molecules can be studied in ion traps due to ionization methods introduced during the past fifteen years which can produce multiply-charged ions from such large molecules. These methods include electrospray ionization (ESI), massive cluster impact ionization, and matrix-assisted laser desorption ionization (MALDI)). ESI and MALDI in particular have become the ionization methods of choice for most large polyatomic molecules such as those mentioned above. In the case of MALDI, singly charged ions usually dominate the population of ions produced. However, in the case of ESI, multiply charged polyatomic molecules usually dominate the population of ions produced. In addition, the population of multiply charged ions produced with ESI has a distribution, or range, of charge states, all of which are substantially greater than +1 or −1. As such, the population of multiply charged ions produced with ESI has a distribution, or range, of mass-to-charge ratios.\nHaving a population of polyatomic molecules present in the chamber of the ion trap which represents a range of mass-to-charge ratios can be a drawback. In particular, the charge state of the polyatomic molecule of interest may be spread out over 10-15 different ionic states which results in a plurality of relatively weak signals when the population of multiply charged polyatomic ions is analyzed. For example, each charge state gives rise to one relatively weak mass spectrum signal when the population of polyatomic ions is subjected to the previously mentioned “mass-selective axial instability mode” of mass spectrometry. Accordingly, there is a need for a method of operating an ion trap which addresses the aforementioned drawback."}
{"text": "Portable electronic devices have gained widespread use and can provide a variety of functions including, for example, telephonic, electronic messaging and other personal information manager (PIM) application functions. Portable electronic devices can include several types of devices including mobile stations such as simple cellular telephones, smart telephones, wireless PDAs, and laptop computers with wireless 802.11 or Bluetooth capabilities. These devices run on a wide variety of networks from data-only networks such as Mobitex and DataTAC to complex voice and data networks such as GSM/GPRS, CDMA, EDGE, UMTS and CDMA2000 networks.\nDevices such as PDAs or smart telephones are generally intended for handheld use and easy portability. Smaller devices are generally desirable for portability. A touchscreen input/output device is particularly useful on such handheld devices as such handheld devices are small and are therefore limited in space available for user input and output devices. Further, the screen content on the touchscreen input/output devices can be modified depending on the functions and operations being performed.\nTouchscreen input/output devices are constructed of a display, such as a liquid crystal display, with a touch-sensitive overlay. These input/output devices suffer from inherent disadvantages relating to user interaction and response, however. In particular, such touchscreen input/output devices fail to provide a user-desirable tactile quality for positively indicating input, providing a poor user-experience.\nImprovements in input/output devices are therefore desirable."}
{"text": "This disclosure relates generally to the field of chemical compound production and more specifically to an improved process and apparatus for using the extremely time-efficient, energy input-efficient, and therefore cost-efficient “collision physics” to achieve chemical compound production by means of the supersonic collision shockwave reaction mechanism.\nFor centuries, the conventional way of chemical reaction was based on mixing two or more chemicals in a stirring, blending or intermixing method with or without the addition of a liquid, often with the assistance of the input of heat, pressure and time, and sometimes further assisted by introducing a catalyst into the reaction. The objective was to cause the original respective bonding of a chemical A and a different chemical B to be loosened, thereby to allow one or more new compounds to be made as a result of a new pairing of A and B. If either or both of A and B are themselves compounds, then a new pairing of a number of variations of the components of A and B."}
{"text": "The present invention relates to laminated or layered assemblies which include a first layer a face of which is secured by adhering to a first face of a second meshed layer. The second layer comprises a base with free loops projecting from a second face of said second layer.\nThe invention applies particularly to such assemblies which are flexible and which are used for making, for example in the form of bands, the female parts of so-called self-gripping fasteners, the loops of which are intended to engage,in a releasable hooting relationship, with complementary male elements, for example hooks, filiform elements with an enlarged head or the like. These female parts of self-gripping fasteners are used, among others, as adjustable and releasable closures for diapers.\nFor cost reasons one is led to manufacture the second layer base in a light and open meshed form, and therefore with a mechanically fragile and dimensionally unstable structure, without however sacrificing the quality of the loops which constitute the critical operative member of the self-gripping fastener female part. The fragility and instability of the base involves processing difficulties, which include, among others, positioning and securing at high speed this female part on its support. There are also utilization difficulties, due to relatively high tearing-out stresses to which the base is submitted and which result from the action of the male elements on the loops.\nOne is therefore led to reinforce the base by adhering it on a more resistant supporting layer, for example a solid, thin and flexible film. This supporting layer can be either an intermediary support, the laminated assembly being then secured to an article, or the article itself.\nHowever some difficulties are met during the manufacturing of such a laminated assembly because the loops have a tendency to also adhere to the supporting layer through the meshes, together with the base material. This destroys their self-gripping ability since the apex of the loops is firmly laid flat against the base and/or the supporting layer and therefore cannot engage the male elements in the desired hooking relationship. As a matter of fact, to effect this lamination, it is generally necessary to firmly press the two assembly members one against another, for example, by means of calenders."}
{"text": "With any internal combustion engine it is desirable that exhaust gases be treated so that they can be safely discharged into the atmosphere. In some engines, particularly of the diesel type, among the most prevalent operating problems is the presence of particulates of varying size, which are carried in the exhaust gas stream.\nThe particulates are normally bits of carbon. They result from the incomplete combustion of hydrocarbon fuels under certain engine operating conditions. However, the operating efficiency of the engine is also a contributing factor to the amount of carbon produced.\nThe presence of relatively large amounts of carbon particles in any exhaust gas stream is evidenced by a dark, smoky, undesirable effluent. Such smoke is not only offensive aesthetically; in large quantities it can be unhealthy. The problem is most evident in the operation of diesel engines.\nMeans have been provided and are known to the prior art, for the elimination or minimization of the particulate content in exhaust discharge streams. It is known, for example, particulates can be eliminated by a suitable smoke filter of proper construction. Eventually, however, the latter can become saturated and/or inoperable, a result that stems from excessive particulate accumulations which block flow passages.\nIt is further known that the overall engine exhaust gas treating process can be expedited to a degree. This is achieved primarily by passing the hot gas stream through a filter medium, a process which is expedited by providing the filter with a catalyst which is capable of promoting combustion of retained particles.\nIt should be appreciated that the generation of carbon particles is prevalent under virtually all diesel engine operating conditions. It is further appreciated that the quality of an exhaust gas stream in any internal combustion engine will vary in accordance with the immediate operating characteristics of the engine.\nFor example, the temperature range experienced by a diesel exhaust gas stream can vary between slightly above ambient air temperature, to temperatures in excess of 1200.degree. F. When the exhaust gas is hot enough, the carbon particles trapped in a filter will, upon contact with the hot gas, be combusted and the filter rejuvenated. However, in diesel powered passager cars, engine operating conditions under which this automatic rejuvenation can occur is seldom reached.\nWhere it is determined that an engine continuously operates under circumstances that the particulates are continuously produced and accumulated in the filter, the particulate retaining bed must be occasionally rejuvenated.\nUnder some circumstances, and as noted above, rejuvenation will consist of merely introducing the hot exhaust gas stream, containing sufficient oxygen, into the filter bed to contact and ignite the retained carbon particles. The combustion of any large and confined carbon accumulation tends to produce temperatures greatly in excess of that of the exhaust gas. The result is that at such excessively high tempeatures, parts of the filter bed or even the filter casing are susceptible to thermal shock, damage or distortion.\nStated otherwise, the filter bed, when formed of metallic fibers or of a similar filter media, is constructed in a manner to be characterized by a sufficient fiber density to remove solid particulate matter. The filter mass, however, should not be so dense as to establish too great a back pressure against the source of the gas.\nTo withstand the high temperatures which can be expected in internal combustion engine operation, the filter media here contemplated is preferably formed of metallic fibers such as stainless steel or the like. Such fibers are capable of being readily bent or deformed. Thus, a mass of randomly disposed fibers can be compressed to a desired density, received within a reaction chamber, and perform the desired particulate removing function. The degree to which it is compressed will determine the size of the gas flow passages therethrough.\nWhen a filter bed of metallic fibers is utilized in an exhaust gas treatment system, the hot gas stream will tend to follow the path of least resistance from the casing inlet to the outlet. As a general rule, for a uniform density filter bed, this path will be through the center of the bed, or along a path determined by the relative positioning of the inlet and outlet.\nAs a consequence, the bed's outer or peripheral edges will receive a minimal amount of retained particulate matter. As a further consequence, the particulate accumulations will be concentrated in those areas or those parts of the bed where gas flow is most concentrated.\nDuring the rejuvenation period when the combustible accumulations are burned off, the thermal intensity of incineration will be concentrated at certain parts of the bed. Thus, while the rejuvenation process is found to be successful, it might have been achieved at the expense of the parts of the bed which have been permanently deformed or otherwise thermally damaged due to prolonged exposure to the excessive heat.\nThis thermal damage is most noticeable at the downstream end of the gas's flow path. Generally, the filter bed will be found to have developed a conically shaped defective area adjacent to the filter casing outlet. It results from the combined effect of the incineration combustion, together with the hot exhaust gas as the latter becomes progressively heated while transversing the filter bed.\nToward overcoming the stated problems endemic to exhaust gas filters, there is presently provided a filter having a bed formed of randomly disposed and compressed metallic fibers. The latter are sufficiently compacted to define a plurality of passages through which the hot exhaust gas will flow.\nThe physical composition of the overall bed is relatively uniform. However, the packed density of the mass is varied in different segments of the bed. Thus, the exhaust gas flow path through the respective bed passages, in following a path of minimal resistance will urge much of the stream toward the filter's peripheral edges.\nIt is therefore an object of the invention to provide an effective exhaust gas filter for removing particulate combustible matter from an exhaust gas stream.\nA further object is to provide an exhaust gas filter which can be safely rejuvenated without jeopardizing the integrity of the filter medium.\nAnother objective is to provide a filter which utilizes a frangible though compressible filter media which has the capability of being compressed to a desired density whereby to regulate the gas flow path through the filter bed."}
{"text": "Drums are the simplest, and most primitive, musical instruments. The drum is a percussive instrument, which produces sound by striking a membrane. The sound is propagated through a membrane, or drumhead, to the drum shell, which is designed to resonate when the membrane is struck. The drumhead is coupled to the drum shell through drum hoops, lugs, and lug or hoop holders. The energy created by striking the drum head is transferred into a wave in the drum shell, producing the distinctive drum sound, a tuned impulse.\nWhile drums usually cannot play different pitches, they are usually tuned. The drum is tuned by tightening or loosening the drumhead by adjusting the lugs and drum hoops. The tighter the drumhead, the higher the pitch propagated by the drumhead.\nMany drummers use a drum kit. Drum kits have several different drums, which can be individually tuned. A drum kit is often composed of various drums, such as a bass or kick drum, snare drums, and tom drums, as well as assorted cymbals and high-hats. When a drummer is drumming, there is substantial vibration throughout the drum kit. Additionally, the various drums can move or flex as they are struck, meaning that the drums, themselves, are vibrating and, therefore, moving. This is especially true of the tom and snare drums.\nA drum mount or support is a sub-classification that includes many different methods of mounting the drum to legs or to other structural elements. The mounts used for drums often degrade the sound, because the drum is held too tightly, damping or attenuating the tuned impulse. Depending on how the drum mount is attached, and how it supports the drum, it can hinder the drum shell resonance, the drumhead, or both. Any drum mount or support that rigidly fixes itself to the drum, whether to the drum shell or the drum hoops, risks damping the sound. A very small number of mounts attach to lugs or lug holders. As currently implemented, these integrated mounting brackets and lugs are a bad idea, because they fail to properly isolate the tuning of the lug from the tuning of the drum mount. The lugs are supposed to be tuning the drum. By adding additional stiffness (force) to the lugs, a fixed drum will change the tuning of the drum. An integrated lug and drum mount is conceptually attractive, because such a device, if it did not adversely affect the tuning of the drum head, would be easier to install on the drum.\nOther mounts attach to the top and/or bottom hoops of the drum shell. Still other supports attach to the drum shell. Rigidly attached supports, whether mounted to the shell or hoops, will damp the vibration of the drum, and may, ultimately, distort the sound through buzzing or rattling, if the support mount is not properly engineered and attached.\nAs a result, a new device is needed, integrating the lug and drum mount, without unduly affecting the tuning of the drum. The new integrated lug and mount should allow the drum shell to resonate freely. Such an integrated lug and mount needs to maintain the drum, in position, without, itself, generating objectionable sounds. The integrated lug and drum mount should be quickly and easily adjustable. It should also allow for quick set-up and break-down of a drum kit. An integrated lug and drum mount that has variable stiffness, and isolates the drum shell from the structural members supporting the drum would be ideal."}
{"text": "The invention relates to a method for transmitting payload information in a radio communication system having at least one radio network controller, at least one base station and at least one subscriber station.\nThe invention also relates to a radio communication system.\nCommunication systems are becoming increasingly important. Significant efforts are being made to link cable-based communication systems to radio communication systems. The resulting hybrid communication systems generate an increase in the number of available services, but also allow greater flexibility in terms of communication. At the same time, devices are being developed which can use various systems (multi-homing).\nConsiderable importance is attached to the radio communication systems due to the subscriber mobility which is made possible.\nIn radio communication systems, information (e.g. voice, image information, video information, SMS [Short Message Service] or other data) is transmitted with the aid of electromagnetic waves via a radio interface between sending and receiving station (base station or subscriber station). In this case, the emission of the electromagnetic waves takes place using carrier frequencies which lie in the frequency band which is provided for the relevant system.\nFrequencies of 900, 1800 and 1900 MHz are used for the GSM (Global System for Mobile Communication) mobile radio system that has been implemented. These systems essentially transfer voice, fax and SMS (Short Message Service) short messages, as well as digital data.\nFrequencies in the frequency band of approx. 2000 MHz are provided for future mobile radio systems using CDMA or TD/CDMA transmission methods, e.g. UMTS (Universal Mobile Telecommunication System) or other third-generation systems. These third-generation systems are developed with the aims of worldwide radio coverage, a large offering of services for data transmission and, most importantly, flexible management of the capacity of the radio interface, which is the interface with the fewest resources in radio communication systems. In the context of these radio communication systems, the flexible management of the radio interface should primarily allow a subscriber station to send and/or receive a large amount of data at high data speed as required.\nThe access of stations to the shared radio resources of the transmission medium, e.g. time, frequency, throughput or space, is governed by multiple access (MA) methods in these radio communication systems.\nIn the case of time division multiple access (TDMA) methods, each send and receive frequency band is divided into time slots, wherein one or more cyclically repeated time slots are assigned to the stations. Using TDMA, the radio resource of time is separated in a station-specific manner.\nIn the case of frequency division multiple access (FDMA) methods, the whole frequency band is divided into narrow-band portions, wherein one or more narrow-band frequency bands are assigned to the stations. Using FDMA, the radio resource of frequency is separated in a station-specific manner.\nIn the case of code division multiple access (CDMA) methods, the throughput/information which has to be transmitted is encoded in a station-specific manner by a scrambling code which includes of a multiplicity of individual so-called chips, whereby the throughput which must be transmitted is scrambled randomly over a wide frequency range in accordance with a code. The scrambling codes which are used by different stations within a cell/base station are mutually orthogonal or essentially orthogonal in each case, whereby a receiver recognizes the signal throughput which is intended for the receiver and suppresses other signals. Using CDMA, the radio resource of throughput is separated in a station-specific manner by scrambling codes.\nIn the case of orthogonal frequency multiple access methods (OFDM), the data is transferred in a broadband manner, wherein the frequency band is divided into equidistant orthogonal subcarriers, such that the simultaneous phase shifting of the subcarriers covers a two-dimensional data flow in the time-frequency range. Using OFDM, the radio resource of frequency is separated in a station-specific manner by orthogonal subcarriers. The combined data symbols which are transferred on the orthogonal subcarriers during a time unit are called OFDM symbols.\nThe multiple access methods can be combined. In this way, a multiplicity of radio communication systems use a combination of the TDMA and FDMA methods, wherein each narrow-band frequency band is divided into time slots.\nFor the purpose of the aforementioned UMTS mobile radio system, a distinction is made between a so-called FDD (frequency division duplex) mode and a TDD (time division duplex) mode. In particular, the TDD mode is characterized in that a shared frequency band is used for the signal transmission in both uplink (UL) direction and in downlink (DL) direction, while the FDD mode uses a different frequency band for the two transmission directions in each case.\nIn radio communication connections of the second and/or third generation, information can be transmitted in a circuit-switched (CS) or packet-switched (PS) manner.\nThe connection between the base station, of which there is at least one, and the subscriber station, of which there is at least one, takes place via a radio communication interface. The base station can support a plurality of radio cells in this case.\nThe base station, of which there is at least one, and a radio network controller (RNC) are usually components of a base station subsystem (RNS radio network subsystem). A radio communication system normally includes a plurality of base station subsystems which are connected to a core network (CN). The radio network controller of the base station subsystem is connected to an access facility (SGSN Serving GPRS Support Node) of the core network in this case.\nIn addition to individual items of payload information, radio communication systems transmit data which is made available to a plurality of users. Such items of payload information include, for example, video streams or other items of broadcast and/or multicast information. The services for transmitting items of payload information which are not merely intended individually for a single subscriber, but are made available to a plurality of subscribers, are grouped under the term MBMS (Multimedia Broadcast/Multicast Service). Different MBMS (Multimedia Broadcast/Multicast Service) services are usually provided as separate data streams by the core network.\nWO 02/19741 A2 discloses a communication system in which an IP multicast is sent to a plurality of devices in an IP network. In order to ensure that, on the network side, there is information about which end devices wish to receive which IP multicast, an “end device membership query” is periodically sent to the end devices from the network side. Upon receipt of a membership query, an end device sends a membership report as a reply to a local router and, in the membership report, identifies those multicast groups which this end device wishes to receive.\n3GPP (3rd Generation Partnership Project) TS 22.146 V5.2.0 Release 5 describes the requirements which must be supported by a 3GPP system (UTRAN [Universal Terrestrial Radio Access Network] and GERAN [GSM EDGE Radio Access Network]), in order to provide broadcast and multicast services (EDGE: Enhanced Data rates for GSM Evolution).\nBefore the items of payload information are made available as a service to a plurality of subscribers, a notification takes place in relation to the subscriber stations of those subscribers who want to use the service, before the actual transmission of the items of payload information of the service. This notification of the subscriber stations is necessary so that the receivers can be configured. Depending on the mode of the subscriber stations (e.g. “connected mode” or “idle mode”), they are notified by a “notification” or “paging”, for example. Group-specific mechanisms are normally used for notification, wherein a plurality of subscriber stations are addressed simultaneously.\nThe transmission of items of broadcast/multicast information as services should take place advantageously. In particular, any avoidable utilization of radio resources should be avoided whenever possible."}
{"text": "1. Technical Field\nThe present invention relates to an image forming apparatus and a process cartridge.\n2. Related Art\nRecently, attention has been focused on increasing the speed and extending the operational lifetime of image forming apparatuses including a charging means, an exposure means, a developing means, a transfer means and a fixing means, in other words, xerographic image forming apparatuses, as a result of technological developments in these members and systems. Similarly, demands for increased response speeds and increased reliability of subsystems have also intensified. In this regard, electrophotographic photoreceptors used for image forming are exposed to large outside electrical and mechanical forces due to chargers, developing devices, transfer devices, cleaners and the like, and thus are susceptible to image defects such as scratches, abrasion, cracking and the like. Therefore, there is specifically a strong demand for improved response speeds and reliability.\nIn order to suppress scratches, abrasions and the like, and to improve operational lifetime, resins having high mechanical strength may be used for electrophotographic photoreceptors."}
{"text": "Currently, there is an unprecedented need for scalable high performance storage and data management, partly due to the wide use of cloud computing. Large distributed storage systems have been developed to satisfy this need for scalable high performance storage and data management. Some of these large distributed storage systems may support a number of tenants, each of which may include one or more clients. Such large distributed storage systems provide isolated storage services to the different tenants. In order to provide these isolated storage services, the distributed storage systems utilize authentication schemes so that each tenant can access the storage services dedicated to that tenant. Some of the distributed storage systems require each tenant to use a provider specific authentication scheme.\nAlthough isolated storage services are desired for most situations involving tenants, there are situations where it may be desirable to share at least some of the storage services for a particular tenant with clients external to that particular tenant. In addition, having the tenants convert to a provider specific authentication scheme is often inconvenient and sometimes impossible.\nThroughout the description, similar reference numbers may be used to identify similar elements."}
{"text": "1. Field of the Invention\n2. Description of the Prior Art\nA conventional container, in consideration of a cost and a purpose, is generally a paper-made or plastics-made container, such as a paper-made or plastics-made cup, bowl, and dish, or a filling container made of this material, such as a packaging container, a heat preservation container, and a lunch box. An existing paper-made container is provided with a single layer that cannot preserve and insulate heat. The container may be provided with two layers with a partition, the cost is, however, extremely high and the cost efficiency is not satisfying since the container is designed for one-time usage.\nFor heat insulation purpose, the existing paper-made cup is held with an additional plastic cup supporter. However, the addition of the plastic cup supporter is in conflict to the original design intention of the one-time-usage paper-made cup since the cup supporter has to be recycled and stored after the paper-made cup is thrown.\nOn the other hand, a heat insulation container composed of foamed material is provided to solve the previous problems. A conventional foamed container is manufactured by the process that the paper sheet is foamed before modeled to a container. However, the foamed paper sheet has a large volume and is difficult to be stacked. Also, a cutting machine is unable to cut the foamed paper sheet into desired shape due to the soft foamed layer. Moreover, stacked foamed paper sheets are unable to be fed into machines due to the large friction of the foamed layer, so that the yield rate is low.\nBesides, some conventional foamed containers have a plurality of deficiencies. Patent publication US 2002/0182347 showed that the foaming layer composed of 60-90% by weight the microcapsules of foaming agent and 10-40% by weight the binder, and that the paper sheet is foamed before shaped into a container so that the problems mentioned above occur. That is, the microcapsules are much more than the binder. It results that the foaming layer may flake off in powder. Also, microcapsules are quite closed to each other and don't have enough space to expand. Moreover, a plurality of adjacent microcapsules may integrate into a larger capsule to form a rough surface. Thus, the process disclosed in US 2002/0182347 is impossible to be introduced in practice.\nAnother application of foamed paper sheet is shown in U.S. Pat. No. 5,498,860. '860 disclosed a card with a braille portion which is selectively protruded. “The binder of the braille portion is softened, and the microcapsules are foamed at the same time”. However, the process is not suitable for manufacturing a container. When the softening point of the binder is the same with the foaming point of the microcapsules, the microcapsules are unable to be rearranged before foaming so that the foamed surface is rough. In addition, a preferred ratio of the binder and the microcapsules for heat insulation is not provided here. More specifically, a foamed braille portion of a card must have different composition from that of a heat insulation container due to the different purposes.\nWO 03/104338 discloses that the foam coating comprising 1-99.9% or preferably 8-80% by weight of binder and preferably 5-25% by weight of the unexpanded. However, the binder is softened and the heat expandable particles are foamed at the same time. WO 03/104338 does not disclose step(s) for redistributing the heat expandable particles which are dried and cannot flow itself under only the effect of gravity without external force."}
{"text": "The present invention comprises methods and compositions for detecting nucleic acid sequences. More particularly, the present invention comprises methods and compositions for detection of specific genetic sequences using nucleic acid target protection strategies. The methods and compositions of the present invention can be used in the detection of microorganisms, for diagnosis of infectious diseases in humans, animals and plants; assays of blood products, and for genetic analysis for use in such areas as early detection of tumors, forensics, paternity determinations, transplantation of tissues or organs and genetic disease determinations.\nMany target and signal amplification methods have been described in the literature, but none are believed to offer the combination of high specificity, simplicity, and speed. General reviews of these methods have been prepared by Landegren, U., et al., Science 242:229-237 (1988) and Lewis, R., Genetic Engineering News 10:1, 54-55 (1990). These methods include polymerase chain reaction (PCR), PCR in situ, ligase amplification reaction (LAR), ligase hybridization, Qxcex2 bacteriophage replicase, transcription-based amplification system (TAS), genomic amplification with transcript sequencing (GAWTS), nucleic acid sequence-based amplification (NASBA) and in situ hybridization. Some of these various techniques are described below.\nPolymerase Chain Reaction (PCR)\nPCR is the nucleic acid amplification method described in U.S. Pat. Nos. 4,683,195 and 4,683,202 to Mullis. PCR consists of repeated cycles of DNA polymerase generated primer extension reactions. The target DNA is heat denatured and two oligonucleotides, which bracket the target sequence on opposite strands of the DNA to be amplified, are hybridized. These oligonucleotides become primers for use with DNA polymerase. The DNA is copied by primer extension to make a second copy of both strands. By repeating the cycle of heat denaturation, primer hybridization and extension, the target DNA can be amplified a million fold or more in about two to four hours. PCR is a molecular biology tool which must be used in conjunction with a detection technique to determine the results of amplification. The advantage of PCR is that it may increase sensitivity by amplifying the amount of target DNA by 1 million to 1 billion fold in approximately 4 hours. The disadvantage is that contamination may cause false positive results, or reduced specificity.\nTranscription-based Amplification System (TAS)\nTAS utilizes RNA transcription to amplify a DNA or RNA target and is described by Kwoh et al. (1989) Proc. Natl. Acad. Sci., USA 86:1173. TAS uses two phases of amplification. In phase 1, a duplex cDNA is formed containing an overhanging, single-stranded T7 transcription promoter by hybridizing a polynucleotide to the target. The DNA is copied by reverse transcriptase into a duplex form. The duplex is heat denatured and a primer is hybridized to the strand opposite that containing the T7 region. Using this primer, reverse transcriptase is again added to create a double stranded cDNA, which now has a double stranded (active) T7 polymerase binding site. T7 RNA polymerase transcribes the duplex to create a large quantity of single-stranded RNA.\nIn phase 2, the primer is hybridized to the new RNA and again converted to duplex cDNA. The duplex is heat denatured and the cycle is continued as before. The advantage of TAS over PCR, in which two copies of the target are generated during each cycle, is that between 10 and 100 copies of each target molecule are produced with each cycle. This means that 106 fold amplification can be achieved in only 4 to 6 cycles. However, this number of amplification cycles requires approximately three to four hours for completion. The major disadvantage of TAS is that it requires numerous steps involving the addition of enzymes and heat denaturation.\nTranscriptions Amplification (3SR)\nIn a modification of TAS, known as 3SR, enzymatic degradation of the RNA of the RNA/DNA heteroduplex is used instead of heat denaturation, as described by Guatelli et al. (1990) Proc. Natl. Acad. Sci. USA 87:1874. RNAse H and all other enzymes are added to the reaction and all steps occur at the same temperature and without further reagent additions. Following this process, amplifications of 106 to 109 have been achieved in one hour at 42xc2x0 C.\nLigation Amplification (LAR/LAS)\nLigation amplification reaction or ligation amplification system uses DNA ligase and four oligonucleotides, two per target strand. This technique is described by Wu, D. Y. and Wallace, R. B. (1989) Genomics 4:560. The oligonucleotides hybridize to adjacent sequences on the target DNA and are joined by the ligase. The reaction is heat denatured and the cycle repeated. LAR suffers from the fact that the ligases can join the oligonucleotides even when they are not hybridized to the target DNA. This results in a high background. In addition, LAR is not an efficient reaction and therefore requires approximately five hours for each cycle. Thus, the amplification requires several days for completion.\nQxcex2 Replicase\nIn this technique, RNA replicase for the bacteriophage Qxcex2, which replicates single-stranded RNA, is used to amplify the target DNA, as described by Lizardi et al. (1988) Bio/Technology 6:1197. First, the target DNA is hybridized to a primer including a T7 promoter and a Qxcex2 5xe2x80x2 sequence region. Using this primer, reverse transcriptase generates a cDNA connecting the primer to its 5xe2x80x2 end in the process. These two steps are similar to the TAS protocol. The resulting heteroduplex is heat denatured. Next, a second primer containing a Qxcex23xe2x80x2 sequence region is used to initiate a second round of cDNA synthesis. This results in a double stranded DNA containing both 5xe2x80x2 and 3xe2x80x2 ends of the Qxcex2 bacteriophage as well as an active T7 RNA polymerase binding site. T7 RNA polymerase then transcribes the double-stranded DNA into new RNA, which mimics the Qxcex2. After extensive washing to remove any unhybridized probe, the new RNA is eluted from the target and replicated by Qxcex2 replicase. The latter reaction creates 107 fold amplification in approximately 20 minutes. Significant background may be formed due to minute amounts of probe RNA that is non-specifically retained during the reaction.\nChiron Signal Amplification\nThe Chiron system, as described by Urdea et al. (1987) Gene 61:253, is extremely complex. It utilizes 12 capture oligonucleotide probes, 36 labeled oligonucleotides, 20 biotinylated immobilization probes that are crosslinked to 20 more enzyme-labeled probes. This massive conglomerate is built-up in a stepwise fashion requiring numerous washing and reagent addition steps. Amplification is limited because there is no cycle. The probes simply form a large network.\nImClone Signal Amplification\nThe ImClone technique utilizes a network concept similar to Chiron, but the approach is completely different. The ImClone technique is described in Kohlbert et al. (1989) Mol. and Cell Probes 3:59. ImClone first binds a single-stranded M13 phage DNA containing targeted probe. To this bound circular DNA is then hybridized about five additional DNA fragments that only bind to one end and the other end hangs freely out in the solution. Another probe set is then hybridized to the hanging portion of the previous set of probes. The latter set is either labeled directly with an enzyme or it is biotinylated. If it is biotinylated, then detection is via a streptavidin enzyme complex. In either case, detection is through an enzyme color reaction. Like the Chiron method, the ImClone method relies on build-up of a large network. Because there is no repeated cycle, the reaction is not geometrically expanded, resulting in limited amplification.\nWhile the nucleic acid amplification methods described above allow for the detection of relatively small quantities of target nucleic acid molecules, there is a need for the ability to detect target nucleic acid molecules in a shorter amount of time with less background interference. Problems inherent in PCR and other amplification techniques involve sample contamination during the collection techniques and the presence of amplicons (amplified target DNA). There are problems with non-specific target amplification mediated by closely related sequences and the production of primer dimers. There is also poor control of specificity, resulting in false positive reactions, and poor control of sensitivity, resulting in false negative reactions. PCR results must often be confirmed and validated by other techniques such as probe hybridization, Southern blotting or in situ hybridization.\nAdditionally, PCR and amplification techniques can only be used with very small amounts of starting sample DNA, in the range of a maximum of 1 microgram. This negates use of PCR techniques for the detection of low copy number nucleic acid targets. For example, early detection of HIV infection, soon after the initial viral infection, would be almost impossible to detect using PCR.\nThus, compositions, methods and kits are needed that are capable of detecting specific nucleic acid sequences and isolating them. Especially needed are methods and kits that would allow for the detection of low copy number nucleic acid target sequences. Additionally, there is need for methods and kits that provide the flexibility that would allow for isolation of nucleic acid sequences using a desired level of specificity.\nWhat is also needed are methods that do not use amplification techniques, but do allow for the isolation of a specific target sequence from any amount of starting nucleic acid, especially large amounts, and have the flexibility to accomplish the isolation at several levels of specificity, depending on the level of specificity desired.\nIn accordance with the present invention, methods and compositions are provided for the detection of specific nucleic acid sequences from cellular or tissue sources. More particularly, the present invention includes methods and compositions for the detection of nucleic acid sequences using a protection molecule that forms a protected nucleic acid sequence (PNAS) such as a triplex or duplex nucleic acid structure that includes the target nucleic acid sequence. The targent nucleic acid sequence is the specific sequence being detected. An assay using the methods of the present invention may include one, two or three levels of specificity to minimize false positive signals. An assay using the methods or compositions of the present invention can be performed on large amounts of purified DNA in a single test, with high levels of sensitivity, thus eliminating the need for in vitro DNA amplification procedures.\nWhen the target nucleic acid sequence is double-stranded, the structure formed with the protection molecule is a triplex. When the target nucleic acid sequence is single-stranded, the structure formed with the protection molecule is a duplex. In this disclosure, where triplex structures are discussed, one can also substitute duplex structures or structures using PNA (peptide-nucleic acid) and the appropriate nucleases. Assays using the methods of the present invention may be referred to as TPA, Target Protection Assays.\nThe initial level of specificity utilizes protection molecules such as oligonucleotides or peptide-nucleic acids (PNA) to bind to specific target sequences of interest. Such binding may be accomplished by formation of Hoogstein-type hydrogen bonds. The protection molecule, bound to the target nucleic acid sequence, forms the protected nucleic acid sequence (PNAS). Once these PNAS structures are formed and stabilized in solution, the non-specific DNA is digested. For example, this digestion can be accomplished with a combination of endonucleases and a double-strand-dependent exonuclease, such as DNA Exonuclease III (Exo III). The endonucleases used in this example are designed to cut on both sides of the PNAS, leaving approximately 20 base pairs of DNA on each side of the sequence. Exo III, an exonuclease which progressively cleaves one strand of the DNA from the 3xe2x80x2 end, is inhibited by the triple helix structure. Using a combination of nucleases, the unprotected DNA sequences are digested completely. A method of the present invention involving a lower level of specificity would employ an affinity molecule for capture and a reporter molecule for labeling in conjunction with the protection probe.\nHowever, if a higher level of specificity is required, 5xe2x80x2 flanking regions can be generated on either or both sides of the PNAS to allow for assays employing two further levels of specificity. The structure formed, a PNAS with flanking regions is termed PNAS/tail. Following the selected digestion around the PNAS, a capture probe, such as an oligonucleotide complementary to one of the single stranded flanking regions, is added. The capture probe is allowed to hybridize to a single-stranded region. For example, the capture probe could be an oligonucleotide that would bind to a single-stranded region and have an affinity molecule attached. For example, the affinity molecule could be didoxigenein or biotin. The capture probe comprises an affinity molecule and is capable of associating with the PNAS.\nA capturing system is used to isolate the PNAS with the capture probe attached. Any capturing system that is capable of binding to the capture probe and separating the PNAS/tails with affinity molecule from the mixture is contemplated. In the example used above, such a capture system may comprise using magnetic beads coated with anti-didoxigenein antibodies for binding to the didoxigenein-capture probe portion or, streptavidin for binding to the biotin-capture probe portion. The PNAS/tails with affinity molecule, now attached to the magnetic beads, are separated from non-specific complexes and washed to remove any non-specific nucleic acid sequences. Such washing may use any washing technique known in the art. For example, a magnetic particle holder could be used. Again, should this be the level of specificity required, the present invention comprises assays that also have a reporter molecule associated with the protection molecule or the capture probe.\nA third level of specific detection involves the addition of a labeled reporter probe. The reporter probe comprises a detectable label and is capable of associating with the PNAS. For example, the reporter probed may comprise an oligonucleotide complementary to the 5xe2x80x2 single-stranded tail that is part of the PNAS/tail. This 5xe2x80x2 region may or may not be on the opposite flanking tail to which the capture probe binds. The reporter probe may be labelled with any labels known in the art such as radioactivity or non-radioactive labels such as labeled with biotin or didoxigenein for indirect detection, or directly with a fluorescent reporter molecule, e.g., fluorescein, or chemiluminescent or bioluminescent labels. An excess of reporter probe is added to the washed magnetic bead-triplex complex and allowed to hybridize. Detection of the bound labelled reporter probe can be accomplished after washing by using detection devices specific for the type of label used. For example, if a fluorescent labeled reporter probe is used, the labeled sequences can be detected using a fluorometer or viewing the beads through a fluorescent microscope. Alternatively, the amount of bound probe can be directly assessed by fluorescent anisotropy with an analyzer such as the Abbott TDM analyzer.\nCompositions of the present invention include compositions comprising the components to practice the methods taught herein. For example, a composition comprising a labeled protection molecule with an affinity molecule could be used in an assay with a first level of specificity. A composition comprising a labeled protection molecule and a capture probe could be used in a level two specificity assay. A composition comprising a protection molecule, a capture probe and a reporter probe could be used in a level three assay. It is to be understood that the individual molecules, probes and components can also be provided individually.\nThe present invention is especially useful for detecting specific genetic sequences. The present invention comprises methods such as the Target Protection Assay (TPA) in all its formats, which have the advantage of allowing the processing of very large amounts of purified nucleic acids, thus eliminating the need for artificial amplification procedures such as PCR, while enabling the detection of a specific target sequence. In addition, the three levels of specificityxe2x80x94PNAS formation, capture probe binding, and reporter probe bindingxe2x80x94reduce technical problems such as those associated with false positive signals from non-specific amplification and/or hybridization.\nThe present invention comprises a method for detecting a target nucleic acid sequence, comprising obtaining isolated nucleic acid sequences from a sample suspected of containing a target nucleic acid sequence; contacting a protection molecule with the nucleic acid sequences under hydridizing conditions sufficient to form a PNAS; and detecting the PNAS. The methods may further comprise the steps of digesting the isolated nucleic acids containing one or more PNAS with nucleolytic enzymes to form a PNAS/tail; and hybridizing a capture molecule to the PNAS/tail; prior to the step of detecting the PNAS. Additionally, the methods may further comprise the step of hybridizing of a reporter molecule to the PNAS/tail; prior to the step of detecting the PNAS. A method for detecting specific nucleic acid sequences, comprising obtaining isolated nucleic acid sequences from a sample suspected of containing a target nucleic acid sequence; contacting a protection molecule with the nucleic acid sequences under hydridizing conditions sufficient to form a PNAS; digesting the isolated nucleic acids containing one or more PNAS with nucleolytic enzymes to form a PNAS/tail; hybridizing a capture molecule to the PNAS/tail; hybridizing of a reporter molecule to the PNAS/tail; and detecting the PNAS.\nThe present invention comprises compositions for detecting specific nucleic acid sequences, comprising a protection molecule capable of binding with a specific nucleic acid sequence. A composition of the present invention may further comprise a capture molecule. Additionally, a composition of the present invention may further comprise a reporter molecule.\nThe methods and compositions of the present invention should be ideal for the detection of viruses and other microorganisms such as pathogens of humans, animals and plants, as well as genetic analysis of polymorphic gene sequences such as HLA typing. The methods of the present invention can be used in forensics, paternity determinations, or transplantation or organs or tissues, or genetic disease analysis.\nAccordingly, it is an object of the present invention to provide methods to detect specific genetic sequences.\nIt is yet another object of the present invention to provide methods for detecting specific DNA sequences involving triplex nucleotide structures.\nIt is another object of the present invention to provide methods for detecting specific RNA sequences involving triplex nucleotide structures.\nIt is yet another object of the present invention to provide methods for detecting specific DNA sequences involving duplex nucleotide structures.\nIt is another object of the present invention to provide methods for detecting specific RNA sequences involving duplex nucleotide structures.\nIt is another object of the present invention to provide methods for detecting specific RNA sequences involving PNA structures.\nIt is another object of the present invention to provide methods for detecting specific DNA sequences involving PNA structures.\nIt is yet another object of the present invention to provide methods for detecting specific DNA sequences involving antibodies.\nIt is yet another object of the present invention to provide methods for detecting specific RNA sequences involving antibodies.\nAnother object of the present invention is to provide a method of detecting nucleic acid sequences involving radioactive labeled nucleic acids.\nIt is another object of the present invention to provide a method of detecting nucleic acid sequences involving non-radioactive labeled nucleic acids.\nIt is yet another object of the present invention to provide a method of detection of specific genetic sequences with variable levels of specificity.\nAnother object of the present invention is to provide a method of detecting nucleic acid sequences for the determination of the identity of microorganisms.\nIt is another object of the present invention to provide a method of detecting nucleic acid sequences for the determination of the identity of human pathogens.\nIt is yet another object of the present invention to provide a method of detecting nucleic acid sequences for the determination of the identity of animal pathogens.\nIt is yet another object of the present invention to provide a method of detecting nucleic acid sequences for the determination of the identity of plant pathogens.\nIt is another object of the present invention to provide a method of detecting nucleic acid sequences for the determination of the genetic relationship, such as paternity or species identification, of a sample.\nIt is yet another object of the present invention to provide a method of detecting nucleic acid sequences for the determination of potential donors of organs or tissues for transplantation purposes or for protecting the blood supply.\nIt is another object of the present invention to provide a method of detecting nucleic acid sequences for use in forensic determinations.\nIt is yet another object of the present invention to provide a method of detecting nucleic acid sequences for the analysis of genetic diseases.\nIt is another object of the present invention to provide methods for testing body or tissue fluids to detect microorganisms or other pathogens.\nThese and other objects, features and advantages of the present invention will become apparent after a review of the following detailed description of the disclosed embodiments and the appended claims."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for processing a digital image and a method of controlling the apparatus, and more particularly, to an apparatus for processing a digital image that can capture an image input through an image pickup device by pressing a shutter-release button, and a method of controlling the apparatus.\n2. Description of the Related Art\nIn general, examples of apparatuses for processing a digital image include all devices that process images or use an image recognition sensor, such as a digital camera, a personal digital assistant (PDA), a phone camera, and a personal computer (PC) camera.\nAn apparatus for processing a digital image, such as a camera, can display images input through an image pickup device. When a user presses a shutter-release button, the apparatus for processing the digital image can capture a desired image and store the captured image as an image file.\nWhen the user presses the shutter-release button, however, the apparatus for processing the digital image may be shaken due to the user's hand shake, resulting in blurry and unclear images and a failure in photographic quality.\nAs the apparatus for processing the digital image has recently become smaller and slimmer, the apparatus for processing the digital image has become more affected by hand shake caused when the shutter-release button is pressed. In addition, as the apparatus for processing the digital image has been developed to offer higher magnification and higher resolution, the risk of shaking the apparatus for processing the digital image due to hand shake caused when the shutter-release button is pressed has increased.\nTo solve the problem that a desired image cannot be captured due to hand shake caused when a user presses a shutter-release button, a hand shake sensor may be used. The hand shake sensor can measure the amount of hand shake and correct for it. However, the hand shake sensor is then additionally required."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a method of and an apparatus for detecting atrial fibrillation by monitoring and analyzing pulse beats using an algorithm that can exclude many rhythm abnormalities that are not atrial fibrillation.\n2. Discussion of the Related Art\nThe heart is the major muscle that functions as the primary pump for blood flow throughout the body. The heart contains two upper chambers called atria and two lower chambers called ventricles. The right atrium receives oxygen-depleted blood while the left atrium receives blood enriched with oxygen from the lungs. When the atria are full, the outlet valves within the heart open and the atria squeeze blood into the ventricles. The right ventricle then pumps oxygen-depleted blood to the lungs while the left ventricle pumps oxygen-enriched blood to all parts of the body. In this fashion, the heart functions primarily as a double sided pump.\nThe heart's internal pacemaker, known as the sinus node, signals the start of each heart beat. This signal originates in the right atrium in the sinoatrial node and travels simultaneously to the left atrium and down to the interatrial septum to the atrioventricular node. The cycle of electrical stimulation that normally occurs is referred to as normal sinus rhythm. The contraction of the ventricles will be referred to as the heart beats.\nMany rhythm abnormalities may be present. Atrial fibrillation is one rhythm abnormality in which the atria do not contract normally. Instead, there is a continuously varying pattern of electrical activation of the atria resulting in a rapid highly irregular pattern of impulses reaching the atrioventricular node. The atrioventricular node acts as a filter and allows a reduced number of these impulses to reach the ventricles which results in a highly irregular heartbeat pattern. This irregular pattern has been shown in previous studies to be a random pattern (Bootsma et al: Analysis of R-R Intervals in Patients with Atrial Fibrillation at Rest and During Exercise. Circulation 41: 783, 1970). Whenever the term “irregular’ is used in this application it refers to this random pattern of beats found almost exclusively in atrial fibrillation.\nAtrial fibrillation is one of the most common arrhythmias requiring medical attention. Atrial fibrillation may be caused by a number of heart conditions, such as angina, myocardial infarction, heart valve abnormalities, and high blood pressure. These conditions may stretch or scar the atria, thereby causing irregularities in the heart system. Atrial fibrillation may also accompany lung problems or thyroid gland disorders and is also associated with significant morbidity and possible mortality. All persons, young and old, female or male, including the visually and/or sight impaired, may experience atrial fibrillation.\nAtrial fibrillation may occur intermittently or chronically. The most serious complication of atrial fibrillation is formation of a blood clot in the left atrium which may result in a stroke. Many people who develop atrial fibrillation, however, are unaware of their abnormal rhythm. Some in the medical profession have, therefore, advocated self screening of the pulse to detect for the possible occurrence of atrial fibrillation. The literature, however, is generally limited to disclosing instructions for manually taking one's pulse accompanied with additional descriptive information.\nThe reason for using the pulse to detect atrial fibrillation is that the pulse usually corresponds to the heartbeat. The contraction of the left ventricle ejects blood from the heart into the aorta and the resulting pressure wave is detected as a pulse in the arterial system. However, when atrial fibrillation is present, the amount of time between beats varies irregularly.\nWith a longer time interval between beats, there is more time to fill the ventricles with blood and more blood is ejected by the ventricle in the beat following this long interval. This larger volume of blood in the aorta results in a higher systolic pressure for that beat.\nConversely, when the time interval between beats is short, there is less time for ventricular filling and the volume of blood ejected in the beat following the short time interval is less. This results in a lower systolic pressure for that beat. In some cases, the time interval between beats is so short that the systolic pressure of the following beat is so low that it cannot be palpated as a pulse. A ventricular contraction that cannot be palpated as a pulse in the arterial system results in what is called a “pulse deficit.” This is very common in atrial fibrillation. This pulse deficit means that an irregular pattern of heartbeats in the ventricle may result in a less irregular pulse beat pattern since the shortest intervals between heartbeats may not be detected in the pulse. Therefore, any method used to determine the presence of atrial fibrillation by analyzing the time intervals between beats in the ventricles may not be valid when applied to the pulse beats.\nAn article by Bert K. Bootsma, Adriann J. Hoelen, Jan Strackee and Frits L. Meijler, entitled Analysis of R-R Intervals in Patients with Atrial Fibrillation at Rest and During Exercise, Circulation, Volume XLI, May 1970 describes an analysis of the time intervals between ventricular contractions using the electrocardiogram. The article evaluates the standard deviation divided by the mean of the time intervals between ventricular beats in normal subjects and in those with atrial fibrillation. The article finds that atrial fibrillation can be accurately differentiated from normal sinus rhythm using this formula. However, this was based on ventricular contractions determined by the electrocardiogram and was not applied to the pulse beat intervals.\nDue to the presence of a pulse deficit in atrial fibrillation, results based on ventricular contractions determined by the electrocardiogram may not apply to time intervals determined from analyzing the pulse. Furthermore, the extent of the pulse deficit depends on the method used to determine the pulse beats. A method which detects only pulse beats with high systolic pressures will detect fewer pulse beats compared to a more sensitive method. The more sensitive techniques may be better for detecting more pulse beats but they may also give more false positive readings.\nFor example, with a photoplethysmograph using a finger probe with a light source and a photoelectric detector, when the sensitivity of the device is increased, the slightest finger movement is detected as a pulse beat. This device at the highest sensitivity setting detects an irregular pulse in those with normal sinus rhythm due to random noise from finger movement. At the highest sensitivity setting, this device would not be useful to detect atrial fibrillation in the home setting. At the lowest sensitivity setting, very few pulse beats would even be detected. Therefore, any device and algorithm which uses the pulse to detect atrial fibrillation must be designed specifically for the purpose of detecting atrial fibrillation.\nOther rhythm abnormalities that are not associated with an increased risk of stroke are found very commonly. For example, premature atrial contractions or premature ventricular contractions are found in many people, even those with no heart conditions. These rhythm abnormalities are intermittent but they are found much more commonly than atrial fibrillation and generally do not require treatment. Therefore, a device that is designed to detect atrial fibrillation should include an algorithm that would not detect premature beats as an irregular rhythm. If a device designed to detect atrial fibrillation does not include an algorithm to reduce the detection of premature beats, the number of false positive readings would be very high and the device would be of very limited value. Those who might use such a device on a regular basis to detect asymptomatic atrial fibrillation would quickly stop using this device if the vast majority of abnormal readings turn out to be due to premature beats and not atrial fibrillation.\nU.S. Pat. No. 6,485,429 describes a method and apparatus that can detect arrhythmia using an oscillometric blood pressure monitor. This patent discloses a method for detecting multiple rhythm abnormalities including supraventricular premature contractions, ventricular premature contractions, atrial fibrillation, paroxysmal supraventricular tachycardia, sinus tachycardia, and ventricular bradycardia. A study performed by the inventor of the patent using a commercially manufactured embodiment of the patent found that the device only detected 66.6% of patients with atrial fibrillation while detecting 85.7% of premature ventricular contractions (Forstner K W, American Society of Hypertension 16th Annual Meeting 2003, page 25). Premature ventricular contractions are found much more commonly than atrial fibrillation. Therefore a device that has a very high rate of detecting premature ventricular contractions and a significantly lower rate of detecting atrial fibrillation would not be a very useful device if the goal is to detect atrial fibrillation exclusively.\nOther patents that can detect rhythm abnormalities would also be subject to very high false positive rates if they are used to detect atrial fibrillation. U.S. Pat. No. 6,095,984 describes an embodiment (section 1-2-1) where a pulse variation over 0.5% would be detected as abnormal. The most common rhythm abnormality, sinus arrhythmia, by definition (Braunwald, E. Heart Disease A textbook of Cardiovascular Medicine 1992, p 674) has a heart rate that varies by more than 10%. U.S. Pat. No. 6,095,984 did not describe methods and apparatus to look for atrial fibrillation and would not be useful for home monitoring of atrial fibrillation, because the algorithm can result in multiple sources of false readings and the apparatus was not specifically set for optimal detection of pulse beats in atrial fibrillation.\nAn algorithm designed to detect atrial fibrillation by the irregularity of the pulse beat intervals should be designed to reduce the effects of premature beats. One method of reducing the effect of premature beats is to limit the number of beats used in determining the irregularity of the pulse intervals. For example, if premature beats occur on average every twenty beats, then limiting the analysis to only ten beats would reduce the likelihood of a premature beat occurring during the period being analyzed. This method was used in a trial and published in Wiesel et al, PACE, 27:639-643 (2004). In this study, the rhythm of 450 outpatients of which 54 were documented to have atrial fibrillation was analyzed. The study used an algorithm that calculated the irregularity index, defined as the standard deviation of the time intervals between beats divided by the mean of the time intervals between beats, and analyzed only the last ten beats recorded by an automatic oscillometric blood pressure monitor. All recordings with an irregularity index greater than 0.06 were considered to be irregular. This study found that 100% of patients with atrial fibrillation were detected while 16% of patients without atrial fibrillation were also found to be irregular (sensitivity for detecting atrial fibrillation 100%, specificity 84%). This is a marked improvement in the rate of detecting atrial fibrillation over the method used by Forstner.\nThough the sensitivity for detecting atrial fibrillation using the irregularity index applied to the last ten beats recorded by an automatic blood pressure monitor is excellent, there is room for improvement in specificity. Premature beats are common rhythm abnormalities that can reduce the specificity of the irregularity index. For example, a ventricular premature beat occurs earlier than the normal beats but results in a pause following the beat. A premature beat usually results in a lower pulse pressure than the normal beats. This premature beat with its lower pulse pressure would often not be recorded by the blood pressure monitor at all. As a result of this premature beat, there would be a pause equal to the time interval between two normal beats. For a patient with a heart rate of 60 beats per minute, the pulse interval would be one second between normal beats. A premature ventricular beat would result in a pulse interval of two seconds as measured by a blood pressure cuff. The irregularity index for nine time intervals between beats of one second and one interval of two seconds due to one premature ventricular beat would be equal to 0.29, well above the threshold value of 0.06. If the premature beat occurs 0.5 seconds following a normal beat and it is detected by the blood pressure monitor, then there will be a short interval followed by a long interval. With eight intervals of one second, one interval of 0.5 seconds and one interval of 1.5 seconds, the irregularity index would be 0.24, still very abnormal.\nU.S. Pat. No. 6,519,490 issued to Joseph Wiesel discloses a method to reduce the impact of premature beats by deleting only those beats that occur at shorter intervals than the mean beat to beat interval. This method was found to improve the specificity of an automatic blood monitor. An algorithm that can recognize when a premature beat occurs and eliminates both the short intervals and the long intervals associated with these beats could further reduce the false positive rate and improve the specificity for the blood pressure monitors used to detect atrial fibrillation. Some patients have very frequent premature beats occurring every other beat. In that case, all the time intervals are either shorter or longer than the mean of the time intervals. If all the beats are eliminated because they fall outside the upper and lower threshold values around the mean time interval then that pattern is not typical of atrial fibrillation and will be considered a regular rhythm. Atrial fibrillation has continuously varying time intervals around the mean time interval such that there are usually at least a few time intervals near the mean time interval that will not be eliminated by this algorithm.\nOccasionally, premature beats will occur frequently with varying intervals. For example, if the normal beats occur at time intervals of one second, there may be premature beats occurring at 0.5 second, 0.75 second and 0.85 second all in the same patient during the period of time recorded by the blood pressure monitor. If we eliminate some of these intervals because they exceed a threshold value, other intervals may still be present that can cause a high irregularity index. Another algorithm can be used to reduce the false positive rate for those patients with frequent premature beats. In atrial fibrillation there are no normal beats and having more than half of the time intervals almost exactly equal in length during the 10 to 40 seconds of a typical blood pressure monitor reading is unlikely. On the other hand, even with frequent premature beats, more than half the beats are usually normal beats. Thus, if we find that most of the time intervals are almost equal during one blood pressure reading, it is much more likely that this is not atrial fibrillation.\nWhat is needed is a home monitoring method and apparatus to detect the possible presence of atrial fibrillation and communicate this condition to the user so that the user is alerted to consult a medical practitioner for further testing and/or treatment.\nWhat is also needed is a method that can differentiate atrial fibrillation from a normal pulse pattern and from common heart rhythm abnormalities that are not of significant risk such as sinus arrhythmia, atrial premature beats and ventricular premature beats.\nWhat is further needed is a method of and an apparatus for detecting irregular pulse rhythms during a time period and storing this information so that comparisons may be made with the pulse rate rhythms at later times.\nWhat is further needed is a noninvasive and relatively simple method and apparatus that monitors pulse rate irregularities to detect atrial fibrillation, and that is suitable for use of all ages, and by the hearing and/or visually impaired and that is relatively easy to use.\nWhat is still further needed is a monitoring method and apparatus that detects the presence of irregular pulse beats and then displays and stores: i) the number of irregular pulse beats during a pre-selected time interval; and ii) the duration of time between beats during selected intervals.\nYet another need is for a monitoring method and apparatus that determine whether or not a pulse beat pattern is irregular based on algorithmic or heuristic operations performed on selected pulse beat data.\nWhat is needed is a method and apparatus for detecting the presence of atrial fibrillation by detecting an irregular pattern of pulses using a sphygmomanometer.\nWhat is needed is a method and apparatus for detecting the presence of atrial fibrillation by detecting an irregular pattern of pulses using a plethysmograph such as finger probe with a light source and photodetector."}
{"text": "1. Field of the Invention\nThis invention relates to novel N-(haloacetyl)N-(N'-methylenepyrrolidonyl)-2-alkoxyanilines which are useful as herbicides, and to intermediates useful in the synthesis thereof.\n2. Description of the Prior Art\nU.S. Pat. Nos. 3,769,301 and 3,907,544 disclose related N-(acyl-tert-amidoalkyl) acetanilides, including N-methylenepyrrolidonyl derivatives; however, these compounds are substituted with 2,6-dialkyl groups only."}
{"text": "The present invention refers to a dishwasher, particularly of the household type, with a washing basin accommodating at least one removable basket that is vertically adjustable to support dishes of various sizes.\nA dishwasher with an adjustable basket is known, for example from DE 3,103,059, in which the basket slides on horizontal guides. The basket is connected to the guides by means of lateral sliding elements, which must be provided in a minimum quantity of four. The basket is connected eccentrically to the sliding elements. By manual rotation of eccentric elements, the vertical position can be varied selectively between a corresponding upper position and a lower position. Because of the eccentric connection between the basket and the lateral sliding elements, the latter are subject to high bending forces every time the basket is moved horizontally along the guides. Therefore, suitable mechanical locking devices are necessary to maintain the basket securely in the two adjustable positions mentioned above. These locking devices, however, are subject to high loads and make the structure of the machine undesirably complicated. In any case, vertical adjustment of the basket requires an undesirably laborious and uncomfortable manual operation.\nFrom DE 4,228,954, a dishwasher is also known in which the lateral sliding elements are connected to the basket in a vertically sliding manner. The position of the basket can be adjusted to discreet levels defined by several coupling sites made on spring elements of the basket. The coupling sites are engaged selectively by corresponding support indentations in the sliding elements. This solution also turns out to be particularly inconvenient to use, since vertical adjustment of the basket requires separate manual operations of disengaging the spring elements and sliding the basket vertically for each side of the basket."}
{"text": "Wavelength-division-multiplexing (WDM) is a key enabling technology in today's high speed digital communication infrastructure that supports vast amount of data transportation which in turn is essential for many data centric informational applications such as, for example, many internet based applications. The transportation of vast amount of data are made via optical digital signals over an extensive optical fiber network across the country or global and such optical signals, riding on different wavelengths (or “colors”), are distributed and/or re-distributed among various parts or branches of the optical network via signal exchanges located at various data centers and other facilities. Sometimes, such signal distribution and/or re-distribution may also involve conversion of signals from optical to electric, and then from electric to optical, in terms of its carrying media.\nWhen optical signals are distributed and/or re-distributed within an optical network, it often involves interconnecting optical signals from one signal handling unit, which may be, for example, an optical signal transponder installed in a shelf hosted by a rack (“bay”), to another signal handling unit which may be located in a same room, in a different room of a same floor, or sometimes in a different floor. FIG. 15 is a demonstrative illustration of an interconnected optical system 900 as is known in the art. System 900 is simplistically illustrated to include a first bay 910 and a second bay 920 located, for example, in a same room of a building. Bay 910 may include, from top to bottom, multiple shelfs such as shelf 911, with each of the shelfs having multiple optical signal transponders.\nAn optical signal coming from an optical signal transponder located in bay 910 may be connected to another optical signal transponder located in bay 920 via a piece of fiber 902, which may have connectors 901 and 903 at its two ends connecting to the signal transponders. Multiple optical signals from bay 910 may need to be connected to multiple destinations in bay 920, and vise versus, using multiple pieces of fibers. Generally the number of fibers needed equals to the number of optical signals being interconnected between the bays, which is demonstratively illustrated in FIG. 15 by a second piece of fiber 904 and the “dots” in between which represents the existence of many more fibers between bay 910 and bay 920.\nWith the ever increasing data rate, in particular rapid deployment of WDM technology, the number of optical signals of different wavelengths that need to be interconnected between different bays, and sometimes between different shelfs in a same bay, has increased dramatically resulting in the explosive use of fibers in signal interconnection. FIG. 16 is an exemplary picture of a traditional optical system 990 where optical signal interconnect among the bays are provided by an enormous stack of individual pieces of fibers 991 which, like cooked spaghetti, are usually “dumped” at the backside of the various bays."}
{"text": "It is generally customary to employ metal catalysts for the crosslinking of CEC coating layers in order to reduce the stoving temperature. In this connection, organotin catalysts have generally been used. The use of dialkyltin oxide is desirable, and the use of dibutyltin oxide (DBTO) is particularly desirable due to its ready availability and favourable price. However, in order that it can successfully be used in CEC coatings, DBTO requires conditioning. This conditioning may consist of a separate synthesis step, i.e. in the chemical conversion of the DBTO into a dibutyltin dicarboxylate for example.\nThus CEC coating media which crosslink by means of blocked isocyanates are known from EP-A-0 509 437. These contain dialkyltin dicarboxylates derived from aromatic carboxylic acids as crosslinking catalysts, in addition to bismuth or zirconium compounds as additional catalysts. Bismuth hydroxide, trioxide, nitrate, benzoate, citrate, oxycarbonate and silicate can be used as the bismuth compounds.\nDE-A-39 40 781 describes a process for producing catalyst pastes containing dibutyltin oxide by the preliminary dispersion of DBTO with organic solvent and 0.05 to 5 moles, per mole of DBTO, of an organic acid which is customary for the neutralisation of cathodically depositable electro-dip coatings, at a water content of up to 5% by weight with respect to the sum of DBTO, solvent, acid and water, and the subsequent dispersion and comminution of the mixture obtained with grinding binder vehicles and water. Preliminary dispersion is preferably effected within a temperature range from 50 to 80.degree. C. During the production of CEC coating media, the catalyst pastes drag in organic acid. However, low acid contents are wanted in CEC coating media, since they promote good wrap-around behaviour and low current consumption during cathodic lacquer deposition."}
{"text": "Currently, most polarization beam splitters are large in volume, and can hardly be used in optical integrated circuits. However, micro devices including the polarization beam splitter can be manufactured on the basis of photonic crystals. Up to now, there are two kinds of polarization beam splitter based on photonic crystals, one of which utilizes a photonic crystal with an S-polarization-wave bandgap and a P-polarization-wave transmission band, or with a P-polarization-wave bandgap and an S-polarization-wave transmission band to realize the polarization splitting of the wave. This kind of polarization beam splitter can only be used as a separate photonic crystal device in conventional optical waveguides, and is difficult to be integrated into other photonic crystal devices, since the transmittance and degree of polarization thereof are poor. The other kind is based on coupled waveguides with different relative coupling lengths for coupling light waves in different polarization states into different waveguides by means of long-distance coupling between waveguides, utilizing the method of periodic coupling and odd-even state alternation between the waveguides. The polarization beam splitters obtained by the two methods mentioned above, still have a relative large volume and a narrow range of operation wavelength, so they are not conducive to practical applications."}
{"text": "Medicinal herbs have been traditionally used for a variety of purposes, particularly to prevent and treat a variety of diseases including cancer. Some of the most important chemotherapeutic agents used in the treatment of breast cancer are derivatives of herbal remedies. These include the taxanes obtained from the Pacific yew tree, Taxus brevifolia reserpine, and vincristine and vinblastine which are obtained from the Rosy periwinkle shrub. The use of botanical medicine is widespread in all regions of the developing world and is rapidly growing in industrialized countries especially among patients diagnosed with life-threatening diseases such as cancer.\nDespite their broad use, there exists insufficient scientific data on the safety and efficacy of many of the herbal therapies. Recently, a large in vitro study on the effect of more than 70 Chinese herbs extracts on growth of five breast cancer cell lines found that 21% of the herbal extracts demonstrated more than 50% growth inhibition on at least 4 from the 5 cell lines [1].\nThe emphasis of the Chinese medicine is on the treatment of the body as a whole. Chinese medicine views include a three-fold approach of using herbs that treat (I) the vitality (tonifying herbs), (II) the cancer (anti-cancer herbs) and (III) the side effects of chemotherapy.\nIn the Materia Medica of Chinese medicine [2], these herbs are selected from 8 classic categories (which are used traditionally to categorize Chinese herbs): (1) herbs that tonify the Qi; (2) herbs that tonify Blood; (3) herbs that clear heat and relieve toxicity; (4) herbs that clear heat and fire; (5) herbs that regulate the flow of Qi; (6) herbs that drain “dampness”; (7) herbs that invigorate the blood; and (8) herbs that tonify the Yin."}
{"text": "Hazard detectors use sensors to detect substances in the air that may be harmful or that may indicate the development of a hazardous situation. For example, carbon monoxide (CO) and radon gas are substances that can be harmful to humans and animals if exposed to high amounts. However, these substances are difficult to detect with the human senses because they are colorless, odorless, and tasteless. A hazard detector can detect the presence of these substances and prevent the harmful effects of exposure by alarming to notify a user. In other instances, a substance such as smoke, while not necessarily harmful in and of itself, can indicate the development of a hazardous situation, such as fire. An early alarm of the presence of such a substance can prevent the hazardous situation from developing or minimize the harmful effects of the situation. Interconnected hazard detectors include detectors that are connected to a network, enabling communication between the detectors or with a central control unit. This provides several advantages over standalone detectors, including the ability to activate multiple alarms when a single detector is triggered. Hazard detectors may be certified under standards defined by governing bodies and/or by companies that perform safety testing, such as Underwriters Laboratories (UL). For example, certain UL standards define thresholds for when smoke detectors and CO detectors should sound an alarm. Certain UL standards also define the required characteristics of the alarm, such as powering requirements and the volume, pitch, and pattern of the alarming sound."}
{"text": "Generating a multimedia document, such as a yearbook or slideshow, that includes a set of images that is representative of different demographic groups (e.g., a group of individuals in their 20s, a group of individuals in their 30s, etc.) within a large group of members (e.g., employees at a company) can be a difficult task. For example, it takes a lot of effort and time in creating an end of the year slideshow for an organization (e.g., company, non-profit organization) that includes pictures of members of the organization (e.g., employees of company). The creator(s) of the slideshow may want to ensure that each demographic group (e.g., groups based on age, race, gender) are not overly or inadequately represented in comparison to the other demographic groups in order to provide fairness and increase member enjoyment. For example, the creator(s) of the yearbook may want to ensure that the slideshow does not include many images of a particular demographic group (e.g., twenty pictures of males) while only including a single image of another class demographic group (e.g., a single picture of a female). However, the creator(s) of the slideshow may have hundreds or thousands of images of hundreds of members of the organization to choose to include in the slideshow. As a result, it may be overwhelming for the creator(s) of the slideshow to ensure that each member of the organization is adequately represented in the slideshow and that there are not some members that are overly represented in the yearbook.\nCurrently, a multimedia document, such as a slideshow, may be generated by selecting an image for inclusion in the slideshow, where the image has associated facial detection information. A face location is determined in the selected image based on the facial detection information and the selected image is cropped based on the determined face location to generate a cropped image depicting the included face. The cropped image is inserted into a slide associated with the slideshow. However, there is not currently an approach in ensuring a specific distribution of images in the slideshow for various demographic groups, such as ensuring a particular demographic group is not overly or insufficiently represented in the slideshow.\nFurthermore, a multimedia document, such as a yearbook, may be designed and created utilizing a facial recognition module for identifying a subject in a photo and tagging the photo with an identity of the subject. Identifying the subject may include comparing the photo to a plurality of previously tagged photos to identify the subject. However, there is not currently an approach in ensuring a specific distribution of images for demographic groups in the yearbook. For example, there is not currently an approach in ensuring that each demographic group has a minimum number of images in the yearbook. Neither is there currently an approach in ensuring that there is a certain percentage of images out of the total number of images in the yearbook for a particular demographic group.\nTherefore, there is a need in the art for ensuring a desired distribution of images for various demographic groups in a multimedia document (e.g., yearbook, slideshow, video, website)."}
{"text": "1. Field of the Invention\nThe present invention relates generally to furniture and, more specifically to interactive furniture for dieters designed to enable an obese person to engage in a technologically induced environ which promotes behavior modification, socialization and knowledge. Obese people often have a multiplicity of problems ranging from physical and social isolation, a sedentary lifestyle and a lack of basic nutritional knowledge. Each problem exacerbates the condition of obesity thereby creating a cycle of despair, defeat, and an increasing sense of isolation thus leading to increased eating and inertia. The present invention seeks to provide technological support through the use of a multi-media computer built into furniture designed exclusively for the modification of negative behaviors contributing to obesity. Additionally it seeks to create an environment conducive to healthy eating habits, exercise, increased self-esteem and a more rewarding lifestyle. The present invention coordinates sensory, dietary and exercise components into one integrated unit thereby associating eating, exercising, socializing and health monitoring within one life promoting environ.\nThe present invention includes a horseshoe-shaped table with computer and dual-purpose seat that adjusts to sizes sufficient to comfortably accommodate a user of various sizes and weights and which serves as both a chair and a seated weight scale that inputs the data on weight into the computer. Integrated within the table is a daily food scale where all of the dieter\"\"s food for the day is placed into and weighed with all data being fed into the computer. A smaller scale is also structured within the table that measures the weight of each individual meal and inputs the data into the computer to calculate the caloric content of that meal. A partitioned plate with removable segments is situated on the meal scale to allow for accurate measurement of different food groups by the computer. It can be removed and washed after use. The table has two personal external sensors. One is an elasticized band with an imbedded sensor which fits around the head and chin and over the mandible or jaw of the user and collects information on the amount of chewing or mastication and will be used to teach the user to slow the eating process. A pulse sensor is worn around the user\"\"s wrist and the data inputted therefrom will be used to mobilize the user into exercising before and after eating a meal. The data gathered by the various sensors and scales will permit the application of the cyber personal trainer to fulfill the functions of a real dietary councilor. The present invention will include several modes of support including on-line support that will alleviate social isolation by encouraging networking with similar people and allow inputted individual data to be collectively shared. People used to eating alone in shame and isolation could eat a cyber-regulated meal together. Instead of on-line chat rooms there could be on-line meal rooms. Another support mode will ensure training with privacy as the user utilizes the cyber personal trainer application off-line. This application will take inputted data on the user and structure a customized diet and exercise program for them while providing actual nutrition lessons.\n2. Description of the Prior Art\nNumerous diet control systems have been provided in prior art. Typical of these is U.S. Pat. No. 4,387,777 issued to Stephen R. Ash on Jun. 14, 1983.\nAnother patent was issued to Yianni Attikiouzel on Mar. 27, 1990 as U.S. Pat. No. 4,911,256. Yet another U.S. Pat. No. 4,924,389 was issued to Claude Gerbaulet et al. on May 8, 1990 and still yet another was issued on Aug. 13, 1993 to Mary J. Kretsch et al. as U.S. Pat. No. 5,233,520. Another patent was issued on May 2, 1995 to Gunes M. Ecer as U.S. Pat. No. 5,412,564.\nCaloric counting method and apparatus comprises a base support structure, a plate support structure pivotally mounted on the base support structure, and scale means connected to said plate support structure at a point remote from the pivot line thereof. A plate carrying a combination of food components constituting an entire meal is positioned on the plate support structure to determine the total caloric content of the entire meal. The plate is provided with indicia defining different zones on the plate for foods having different calories per gram, these zones being arranged so that the zone corresponding to foods having the least calories per gram is positioned nearest the pivotal mounting of the plate support structure when the plate is positioned thereon. The patient places the food components constituting an entire meal in the correct zones on the plate and the total caloric content of the entire meal is determined in a single operation by reading the scale which is calibrated in calories.\nA dietetic measurement apparatus having computer means having a controlling instruction means to provide a determination of basic integers for a mass of a food item weighed on a weighing means. The weighing means provides a weighing signal, in the form of a logical periodic signal, to a central processor unit of the computer means, which measures the frequency of the weighing signal and correlates the mass of the food item and the quantities of the nutrients. The nutrients comprise carbohydrate, calcium, cholesterol, fat, fiber, iron, protein, sodium and calorific content. Also a number of command keys are provided and include a ZERO key and a BOWL key to aid in weighing procedures.\nA first memory stores a calorie goal and quantities of food in several categories which a person anticipates will be consumed during a predetermined day. The quantity of food eaten during the preceding days are stored in a second memory. A computer determines the recommended number of calories and quantities of food in each category for the predetermined day as a function of the quantities consumed during the preceding days. The recommended quantities are stored in a third memory and comparisons are made between the estimated quantities and the recommended quantities, the results of which are displayed. The device, used with or without scales, provides for a healthy and balanced nutrition.\nAn interactive computerized dietary measurement system and process which can be used by lay people for accurate measurement of the intake of foods, nutrients, and other food components in the diet. The system includes a computer device coupled with an electronic balance, output/display device(s), user input element(s), a food codes database and a storage element. The system has the ability to provide a plurality of weigh-in measurement options with respect to different food service settings or habits and to respectively tag the measurement of each of the options selected by the user in the storage element so as to track the measurements. The system can instruct user on the next measurement action and required input so that users are not required to remember all of the measurement actions and sequences of the actions for the measurements. The system also provides a plurality of weigh-out measurement options in which the computer will track the measurement record and remind the user of what has been done and what should be done so that a complete and accurate collection of dietary intake data can be achieved. The system is further capable of signaling to the user any weighing errors for allowing user to correct the errors, and of permitting user to make comments or notes as to measurements.\nA system and a method for recording and monitoring of dietary consumption by a consumer is disclosed. The system consists of a computer for storing and processing nutritional information of the type used for diet control, a real time clock for maintaining current date record, a product code entry terminal or a bar code reader for inputting product identification information, a read-write unit adapted to receive one or more integrated circuit (IC) cards of the smart card type having memory and a microprocessor for reading and writing nutritional information into and from a smart card, a printer for printing nutritional information, and optionally an electric display for displaying such information. Consumer inserts his/her personalized smart card into card reader writer before a purchase transaction starts at a food store or a restaurant checkout counter. As foods and drinks are purchased, dietary nutritional consumption data is electronically collected, sorted, and combined with historical daily averages of nutritional consumption data stored in consumer\"\"s personalized smart card and daily averages of nutrition consumption along with other nutritional and personal data is printed out as a report to the consumer at the end of the purchase transaction. Nutritional data may include personal consumption data on calories, calories from fat, calories from sugars, and cholesterol, dietary fibers, sodium, carbohydrates, proteins, and the like. Consumer\"\"s smart card serves as a personal nutritional consumption history file, and is updated each time it is used.\nWhile these diet control systems may be suitable for the purposes for which they were designed, they would not be as suitable for the purposes of the present invention, as hereinafter described.\nThe present invention discloses a horseshoe-shaped table having a computer and a dual-purpose chair associated therewith. The chair adjusts to variously sized users by using a piston and piston support mechanism disposed internally of the chair. The chair also serves as a weight scale or sensor which is disposed in the seat thereof and has data input means, including a headset with microphone and chewing sensor disposed thereon. The arm of the chair has a blood pressure and pulse sensor disposed thereon. Data gathered from the sensor input devices is transmitted to a computer disposed internally of the table having a monitor and keypad disposed on the tabletop. Also shown are a pair of scales which measure both the individual meal food and the daily food allowance which data is also input into the computer. The scales are equipped with segmented platters and non-segmented platters, having various food compartments for receiving different types of food. The monitor is disposed on a suspension arm which is adjustable and is provided with a control crank for adjusting the arm. The data inputs are disposed near the computer in order to receive input data from the input devices. An alternative embodiment of the present invention is disclosed wherein data is transmitted from the chair via electromagnetic wave means to the table wherein the chair is electrically operated and the monitor, and food scales are mounted internal of the table being covered by panels in the closed position and having a control panel for electrically raising these elements from internal of the tables.\nA primary object of the present invention is to provide interactive furniture for dieter\"\"s that will promote behavior modification, socialization, and knowledge for obese people.\nAn additional object of the present invention is to provide interactive furniture for dieter\"\"s including a table with a computer and two scales and an expandable seat with a plurality of sensors to monitor the user\"\"s weight, vital signs, and mastication rate and input the gathered data into the computer.\nA further object of the present invention is to provide interactive furniture for dieter\"\"s wherein the table scales are used to measure the weight of one days food and the weight and caloric content of each individual meal with the data from each scale then inputted into the computer.\nA yet further object of the present invention is to provide interactive furniture for dieter\"\"s wherein the computer uses all the inputted in data from the chair sensors, the scales and exercise equipment to determine a diet and exercise regimen appropriate for that individual thereby acting as a personal trainer.\nAnother object of the present invention is to provide interactive furniture for dieters wherein the computer allows the user to go on-line to form support groups and exchange data with other users of the present invention.\nAnother object of the present invention is to provide interactive furniture for dieters that is economical in cost to manufacture.\nAnother object of the present invention is to provide interactive furniture having a pulse sensor thereon which will be used to encourage the user to increase the pulse rate by activity or exercise before sitting down to eat.\nAnother object of the present invention is to provide interactive furniture wherein exercise equipment is integrated with the computer as well in order to accept inputed data including distance, coloric content, etc.\nFurther objects of the present invention will appear as the description proceeds.\nTo the accomplishments of the above and related objects, this invention may be embodied in the form illustrated in the accompanying drawings, attention being called to the fact, however, that the drawings are illustrative only, and that changes may be made in the specific construction illustrated and described in the appended claims."}
{"text": "Many miniature electro-explosive devices (EEDs) are fired (activated) by the rapid discharge of the energy stored in a storage capacitor into the bridgewire of the EED by the closure of a switch, typically a Field Effect Transistor (FET). In many designs, the capacitor is not able to be fully discharged because as the capacitor voltage drops (as the discharge progresses), the gate-to-source voltage across the FET drops below the turn-on threshold, essentially turning off the switch. Fully discharging the capacitor is desirable to increase the reliability of firing the EED, but a method is needed to achieve this in a low-cost fashion utilizing a minimum number of components, particularly in applications where the volume available for the firing circuit is limited.\nWhat is needed, therefore, are techniques for fully discharging a storage capacitor."}
{"text": "Teaching devices for use in teaching and learning basic mathematics are known, and include devices using marbles to represent numbers. For example, U.S. Pat. Nos. 2,875,529 and 3,400,472, incorporated herein by reference, disclose marble-based teaching devices with relatively limited capabilities. The teaching device disclosed in the '472 patent is designed only for use in teaching addition. The device in the '529 patent is described for use in teaching addition, subtraction, multiplication, and division. However, the '529 device is relatively constrained in its application, and limits the tactile and aural involvement of the student in the learning process. More specifically, the flow of marbles in the '529 device is caused by the gentle tilting of the device, and the marbles are retained in partially enclosed channels. The gentle slope limits the impact-generated noise as the marbles move from one position to the next. The enclosed channels limit the ability of the student or teacher to touch or otherwise interact with the marbles.\nAccordingly, there is a need for a computational teaching device that can be used to teach and learn counting and all four of the basic mathematical operations, and that will provide the opportunity for significant student involvement in the problem-solving process. There also is a need for a device that easily will demonstrate the connections between basic mathematical problems and patterns, fractions, different number bases, graphing, statistics, and negative numbers."}
{"text": "This invention relates to a method for manufacturing a plywood using an improved veneer bonding technique.\nIn prior art techniques of bonding veneers for plywood manufacture, it is well known that poor adhesion such as puncture and deficient bond strength largely depends upon the extent of drying of veneers, that is, the moisture content of veneers.\nHowever, it is impossible at the present state of the art to dry a number of veneers to an equal moisture content, and veneers are dried to a more or less varying extent. Particularly, the recent deterioration of log stock results in a mixed supply of difficult- and easy-to-dry woods. With the increased extent of drying, over-drying results in excessive shrinkage and corrugation of veneers. With the reduced extent of drying, on the other hand, under-drying results in poor adhesion. In connection with veneer drying, a number of problems arise at the site of manufacture.\nThe only one solution for the present status depends on the development of a practical gluing technique which is completely or substantially independent from variable moisture contents. To essentially eliminate these problems, it is necessary to establish a technique capable of gluing veneers having a high moisture content at low cost.\nOne technique of gluing veneers having a high moisture content is disclosed in Japanese Patent Application Publication No. 54-3929 titled \"Plywood Manufacturing Method\". According to the method of this patent publication, a plywood adhesive fluid of well-known formulation is applied to a veneer and dried thereon, and another veneer is placed on the adhesive-applied veneer followed by hot pressing. Although the results described in the patent publication are satifactory, I found that this technique has the following problems. With a dryer equipped for the purpose of drying the adhesive applied and set to proper drying conditions including temperature and time, the adhesive applied on those veneers having a lower moisture content is overdried and the adhesive applied on those veneers having a higher moisture content is under-dried, failing to dry the adhesive applied to a desired extent for all veneers. As a result of the varying dryness of the adhesive, the bond strength achieved by hot pressing also varies over a wide range.\nIn the examples described in the patent publication, drying conditions are changed in accordance with the moisture content of veneers to be handled, and more specifically, the higher the moisture content of veneers, the higher the drying conditions are set. The dryness of the adhesive applied not only depends on the selected adhesive drying conditions, but largely depends on the moisture content of veneers. The varying bond strength is attributable to this fact. Consequently, the dryer for drying the adhesive applied to veneers must be of size and capacity as large as the conventional veneer dryers. Besides, it is well known that there is a significant difference between sapwood and heartwood portions of the same log and between spaced portions of the same veneer. From a standpoint of view of manufacturing plywoods of consistent quality from a variety of starting logs, the above-mentioned shortcomings of the prior art techniques are difficult to overcome in the actual forms of plywood manufacture requiring large-scale production and prevent commercial application of such techniques.\nAn object of this invention is to eliminate the above-mentioned shortcomings involved in the prior art techniques and to provide a novel and improved technique for firmly bonding veneers, even highly wet veneers if necessary, at low cost independent of the moisture content of veneers so that the technique may be readily applied to commercial manufacture."}
{"text": "Anethole (1-methoxy-4-(1-propenyl)benzene) is of great commercial interest as a fragrance and flavouring because of its characteristic aniseed odor. More particularly, anethole is used as fragrance in washing and cleaning compositions or cosmetics, and as flavoring in the foods industry.\nThe preparation of anethole has long been known in the prior art. For instance, one way of obtaining anethole is from natural sources, for example from fennel oil or anise oil—see RU2064920 and CN102491884.\nHowever, the production of fragrances from natural sources is usually costly, and the amounts thus obtainable are limited. Moreover, the purity or production volume of these fragrances often varies because of variable environmental conditions in the production of the raw materials from which they are isolated. There is therefore a great interest in at least partly replacing fragrances from natural sources with substances obtainable by synthetic means.\nThe prior art discloses various synthetic processes for preparation of anethole. For example, Bauer et al., Common Fragrance and Flavor Materials, 2001, 4th Edition, Wiley-VHC, describes the preparation of anethole by a base-catalyzed rearrangement of 1-methoxy-4-allylbenzene (estragiol).\nA further method of preparing anethole, as described, for example, in SU261380 and SU355144, includes a Friedel-Crafts acylation of anisole with propionyl halides or propionic anhydride, followed by the reduction of the carbonyl group and subsequent elimination of water. The acylation of anisole with propionic anhydride can be effected, inter alia, with ZnCl2 and FeCl3 (see Maslozhirovaya Promyshlennost (1974), volume 9, pages 29-30).\nCN103058835 describes a process for synthesizing anethole, in which a Friedel-Crafts reaction proceeding from anisole and propionyl chloride is conducted, followed by a reduction of the carbonyl group to the corresponding alcohol with the aid of NaBH4 and subsequent elimination of water with p-TsOH and KHSO4.\nDE 2418974 B1 describes a process for preparing anethole, in which, in a first step, anisole is condensed with propionaldehyde in the presence of acidic catalysts to give a mixture of bis(methoxyphenyl)propanes. In a second step, the condensation products are cleaved into p-anethole, or o-anethole and anisole, by using catalytic amounts of acid and by heating to temperatures of 100 to 300° C. in the liquid phase. Specifically, the cleavage of the bis(methoxyphenyl)propanes is conducted in the presence of a catalytic amount of concentrated phosphoric acid at 200° C., with distillative removal of the cleavage products formed at about 5 to 30 Torr. Pure p-anethole is obtained from the distillate by fractional distillation. A disadvantage is that the bis(methoxyphenyl)-propanes are not completely cleaved, i.e. are only partially converted. Moreover, only moderate yields of anetholes are achieved. The long residence time of the bis-(methoxyphenyl)propanes and the cleavage products at elevated temperature additionally leads to increased formation of by-products, unwanted isomers, and oligomers and polymers."}
{"text": "1. Field of the Invention\nThe invention relates to a method for automatically loading a feed line with bulk material, in particular a plastic granular material, in a suction delivery system having at least one material separator, which is provided at a distance to a reservoir. Furthermore, the invention relates to a device for implementation of said method.\n2. Background Description\nPneumatic delivery systems are operated generally in such a way that the setting of the material to air mixture is fixed at the material feeding point. Among experts, this is the load of the feed line. Basically, different designs of material feeding points are known, whereby the problem arises for all of them that an optimal mixture cannot be achieved for varying delivery lengths as a result of the fixed setting.\nPneumatic delivery systems for bulk material are planned and designed in such a way that flexibility is achieved for material distribution. Different delivery targets are served in sequence from a single material source, such as, for example, the reservoir. Based on the set values fixed at the material feeding point, a load results for the delivery in accordance with the fan characteristic and the length up to the material separator. For practical reasons, this is set in such a way that delivery is ensured for common material types and the longest delivery distance. However, the load is smaller if shorter delivery distances are served. Reduced load increases the wear of the feed lines and the receiving devices, for example the material separators, at the delivery target due to higher bulk speeds.\nBy reducing the load, the flow of the bulk material is reduced at the same time, although the bulk material is transported faster through the pipeline. Wearing of the material itself is also given, which creates deposits in the feed line or becomes effective as dust.\nAn apparatus for conveying bulk material, in which a variable amount of a conveying gas may be introduced to the bulk material to be conveyed by a vacuum conveying system, is thus known from DE 100 39 564 A1. However, practical tests have shown that the control of the conveying gas alone does not result in a prevention of the disadvantages mentioned above."}
{"text": "Quantum computers can perform computational tasks by executing quantum algorithms. In some quantum computing architectures, quantum algorithms are executed on superconducting qubits that are subject to noise. The superconducting qubits can be implemented, for example, using circuits that include Josephson-junctions."}
{"text": "Firewalls remain the frontier defense for securing networks that are vital to private industry, government agencies, and the military. They enforce a security policy by inspecting and filtering traffic arriving at or departing from a secure network [3, 21, 22]. Meaningful inspections often involve a complex process, and if firewalls are to remain effective, they must adapt to constantly changing network demands, technology, and security threats.\nAs firewall technology evolves to meet new demands and threats, it must continue to act transparently to legitimate users, with little or no effect on the perceived network performance. The networks between two communicating machines should be invisible, especially if traffic requires specific network Quality of Service (QoS), such as bounds on the packet delay, jitter, and throughput. The firewall should process the legitimate traffic quickly, efficiently, and never violate the desired QoS. Unfortunately, the firewall can easily become a bottleneck, given increasing traffic loads and network speeds [12, 14, 18, 20, 21]. Packets must be inspected and compared against complex rule sets and tables, which is a time-consuming process. In addition, audit files must be updated with current connection information. As a result, current firewalls have difficulty maintaining QoS guarantees. Since a firewall can be overwhelmed with traffic, it is also susceptible to Denial of Service (DoS) attacks [21]. Such attacks merely overload/saturate the firewall with illegitimate traffic. Audit files and state tables quickly fill available storage space and legitimate traffic queued behind suspicious traffic encounters long delays. Legitimate users notice poor network performance, and, in the worst case, the secure network is disconnected from the outside world.\nInspecting traffic sent between networks, a firewall provides access control, auditing, and traffic control based on a security policy [3, 21, 22]. The security policy is a list of ordered rules, as seen in Table 1, that defines the action to perform on matching packets. A rule r can be viewed as a k-tuple (r[1], r[2], r[3], . . . , r[k]), where r[i] is specified as a variable length prefix.\nTABLE 1Example Security Policy Consisting of Multiple Ordered RulesSourceDestinationNo.Proto.IPPortIPPortAction1TCP140.*** 80accept2TCP150.**120 80accept3TCP140.**130 20accept4UDP150.***3030accept5*****deny\nIn the examples described herein, rules will initially be represented as the 5-tuple (protocol, IP source address, source port number, IP destination address, destination port number). Fields can be fully specified or contain wildcards ‘*’ in standard prefix format. For example the prefix 192.* would represent any IP address that has 192 as the first dotted-decimal number. In addition to the prefixes, each filter rule has an action, which is to accept or deny. An accept action passes the packet into or from the secure network, while deny causes the packet to be discarded. Rules are applied in order to every arriving packet until a match is found; otherwise, a default action is performed [21, 22].\nUnfortunately, security policies may contain anomalies, such as a packet that matches multiple rules, where the matching rules specify conflicting actions. Security policy anomaly detection and correction is the subject of continued research [13, 21] and is not the focus of the present subject matter. Given the type of security policy required, rule sets can become quite large and complex. Security can be further enhanced with connection state and packet audit information. For example, a table can be used to record the state of each connection, which is useful for preventing certain types of attacks (e.g., TCP SYN flood) [21, 22].\nTraditional firewall implementations consist of a single dedicated machine, similar to a router, that sequentially applies the rule set to each arriving packet. However, packet filtering represents a significantly higher processing load than routing decisions [15, 18, 22]. For example, a firewall that interconnects two 100 Mbps networks would have to process over 300,000 packets per second [21]. Successfully handling this high traffic becomes more difficult as rule sets become more complex [4, 14, 22]. Furthermore, firewalls must be capable of processing even more packets as interface speeds increase. In a high-speed environment (e.g. Gigabit Ethernet), a single firewall can easily become a bottleneck and is susceptible to DoS attacks [4, 7, 10, 12]. An attacker could simply inundate the firewall with traffic, delaying or preventing legitimate packets from being processed. Building a faster single firewall is possible [8, 15, 18, 19, 20]; however, the benefits are temporary (traffic loads and interface speeds are increasing); it is not scalable; it is a single point of failure; and it is generally not cost-effective for all installations.\nA data-parallel firewall (or firewall sandwich) where each firewall node implements an entire firewall rule set is another approach to increase the speed of processing traffic [4, 12, 14, 22]. As shown in FIG. 1, the system consists of multiple identical firewalls 100 connected in parallel. Each machine implements the complete security policy and arriving packets are distributed across the machines for processing in parallel [4]. A packet distributor 102 implements a load balancing algorithm to distribute different packets to different firewall nodes. How the load-balancing algorithm distributes packets is vital to the system and is typically implemented as a high-speed switch in commercial products [11, 12]. Although data-parallel firewalls achieve higher throughput than traditional firewalls [4] and have a redundant design, they have difficulty maintaining QoS across networks. For example, legitimate traffic can encounter delays if it is queued behind traffic that requires more processing. Under these circumstances, users notice poor network performance, which is growing concern as more network applications require QoS assurances. Furthermore, stateful inspection requires all traffic from a certain connection or exchange to traverse the same parallel machine, which is difficult to perform at high speeds [14]. Therefore, new firewall architectures are needed to meet the demands of future networks and the challenges of increasing security threats."}
{"text": "Solar energy appears to be the single most attractive source of energy. Solar systems are not liable to wind damage, as is the case with wind energy systems except in extreme weather conditions, which are exceptional, and would damage even conventional energy transmission systems. Solar collector can be installed almost anywhere, unlike tidal energy systems which must be located in coastal areas. Solar energy systems consume no fuel and the energy is, therefore, \"free\", in a sense. It causes no emissions to pollute the environment. It does not appear to require complex transmission and delivery systems such as electrical power lines, gas or oil pipe lines, or other delivery systems, provided it can be collected \"on site\", i.e., where it is to be used.\nSolar energy has however one major disadvantage, namely, that it is available only intermittently, in the hours of daylight. Accordingly, it is obvious that sufficient solar energy must be collected during the hours of daylight, that a sufficient quantity of energy can be stored, for reuse during the hours of darkness when solar energy is not available. Clearly, systems could be engineered to carry out these functions on a very large scale. However, that would then require the collection and storage of large quantities of solar energy which would then have to be transmitted from place to place, for example, to a large number of homes, connected to a central solar collector. This would result in excessive capital costs, and would likely result in relatively high energy transmission losses. Clearly, such centralized systems would fail to take advantage of the principal significant fact concerning solar energy, namely, in that it is substantially universally available over a large portion of the face of the earth, and thus does not require special transmission systems, to transmit the energy from one place to another. The theoretically ideal solar energy system would be self-contained, so that individual systems could be provided for individual buildings or homes, so that each building or home could be essentially energy self-sufficient. In this way, complex transmission or storage systems would not be required.\nIt is however apparent that in the design of any solar energy collector system it will be necessary to collect solar energy falling on a relatively substantial surface area, to provide sufficient energy even for a single home or building. In order to achieve this, numerous proposals have been made for providing solar collectors based on some form of concave curved mirror usually of parabolic shape having a relatively large surface area. The mirror catches the sun's rays falling on all points of the mirror, and concentrates them into a single point, resulting in a high energy concentration of light at that point. Usually, at the focal point of the mirror, a secondary reflector is provided for reflecting the concentrated point of light, and redirecting it, so as to convert the energy of the concentrated light beam into some other form of energy. Typical proposals of this general type are shown in, for example, U.S. Pat. Nos. 4,286,581, 4,295,462, 4,340,031, and 4,696,285. It is, of course, implicit in any of these solar collector devices that the actual mirror itself must be located in the open air. In all of these proposals, the mirror concentrates the solar energy falling on the mirror at the focal point of the mirror. The concentrated light is re-directed into a portion of the collector where the energy of the light beam is then converted into heat. For example, in U.S. Pat. No. 4,286,581, the solar energy is utilized to heat a tubular chamber, or cavity, and a fluid medium is passed through the chamber or cavity, and absorbs the heat from the concentrated light rays. The energy in the fluid is then transported to some other location for use. In these proposals the mirror was usually mounted on some form of a suitable equatorial tracking mount, so that the mirror could always be pointed directly at the sun. Accordingly, in the proposals of the type described above, it was usually necessary for the entire apparatus consisting of the mirror, the secondary reflector, and also the heat conversion chamber, all to be mounted together and moved simultaneously to track the path of the sun. Such systems therefore require highly complex suspension and tracking mechanisms, to provide the facility for accurate tracking of the sun, by all of the components simultaneously and also require heat transmission systems for transmitting the heat from the outside to the inside of a building.\nOther proposals for collecting solar energy have been made of a more simplistic nature, for example as shown in U.S. Pat. No. 4,059,226. In this type of proposal, large chambers are filled with any suitable heat storage medium. In this case, the medium is simply rock or stone. The chambers are fixed in position, and have glass panels through which the sun's rays may pass and heat up the stones. In the particular example, of this patent, insulating panels are adapted to be swung down over the glass panels, when the heat generated by the sun's rays falls below a certain preset minimum temperature.\nSystems of this kind are relatively massive installations, and are suitable in only relatively limited \"sun belt\" portions of the globe, and must obviously be located for optimum operation in a fixed location, i.e., facing south in the northern hemisphere.\nAnother proposal is shown in U.S. Pat. No. 3,988,166. In this patent, a concave mirror is provided, mounted on a suitable tracking mount, and having a secondary light concentration mirror at the focal point. The light falling from the main mirror on the secondary mirror is concentrated into a light beam. A central opening in the main mirror allows this concentrated light beam to pass through the main mirror. An energy collection chamber is provided immediately behind the main mirror, for receiving energy from the concentrated light beam. The converted energy is then transmitted from this chamber, to a location where it may be of use. While this system is somewhat simpler to construct, it still has the disadvantage that the energy conversion takes place out of doors, and requires the transmission of an energy conversion medium, typically, for example, a fluid medium, from the solar collector itself into the building or facility where the energy is required, typically to heat the building or to provide hot water for a domestic hot water system. Thus thick insulation must be provided to limit heat losses in transmission.\nIn a further proposal, in U.S. Pat. 4,373,514 granted on Feb. 15, 1983 to Lois, there is shown a form of solar collector system in which the collimate beam is routed through a conduit, around various angles. This system uses at least three separate mirrors, to bring the light beam along the angled conduit, to a point where it could be converted. However Lois does not disclose an equatorial mount for tracking the sun's path each day, and each day of the year, and consequently would have very little if any efficiency. In addition, the use of three separate mirrors within the conduits produces further heat transmission losses, which accumulate from mirror to mirror. Consequently, the system would be impractical for the purposes of for example domestic heating or heating of individual buildings, and would be extremely expensive."}
{"text": "The demagnetization or removal of a magnetic field may be accomplished in several ways as described at http://www.ndt-ed.org/EducationResources/CommunityCollege/MagParticle/Physics/Demagnetization.htm, on Oct. 12, 2012, which is incorporated by reference herein. One demagnetization approach is to heat a material above its Curie temperature to produce a random orientation of the magnetic domains, which demagnetizes the material. Another demagnetization approach is to subject the material to a reversing and decreasing magnetic field produced by driving a (de)magnetizer with a decreasing alternating current. This AC demagnetization process, shown in FIG. 1 (PRIOR ART), can be accomplished by pulling a component out and away from a coil with AC passing through it. The same can also be accomplished using an electromagnetic yoke with AC selected. Also, many stationary magnetic particle inspection units come with a demagnetization feature that slowly reduces the AC in a coil in which the component is placed. As can be seen in FIG. 1 (PRIOR ART), which depicts a demagnetization hysteresis curve 102, the current passing through a magnetizing coil decreases in accordance with an alternating current having a current curve 104. The demagnetizing field of the magnetizing coil corresponds to a flux curve 106 that corresponds to the current curve 104, where the alternating polarity H field that is produced by the coil results in a smaller and smaller B field being present in the material inside the coil. An alternative demagnetization approach is the subject of the present invention."}
{"text": "There have been an increasing number of machines activated by cordless power supply utilizing electromagnetic inductance, such as electric toothbrushes, cordless telephones, and portable devices (e.g., PTL 1). There have also been developments of machines activated by cordless power supply utilizing magnetic resonance, in relation to wall-hang television sets and personal computers (e.g., PTL 2). To add this, there have been many developments and suggestions for magnetic elements for wireless power transmission capable of feeding high power with high power transmission efficiency, in the field of wireless power transmission.\nHowever, feeding high power with high power transmission efficiency in such a magnetic element for wireless power transmission causes an excessive heat generation. This generated heat in the magnetic element for wireless power transmission may cause a problem to the magnetic element itself or to the functional parts of the element.\nTo address this problem of heat generation, for example, one approach is to adopt a structure of substrate having an in-built coil such as the one disclosed in PTL 3, in which a planar coil conductor to serve as the heat source is covered with a magnetic layer so as to radiate the heat outside through a conductor layer for heat transfer which is provided to the magnetic layer. This, without a doubt, enables radiation of heat generated by the planar coil to the outside."}
{"text": "A commonly assigned patent of G. and S. Olive, U.S. Pat. No. 4,179,462, issued Dec. 18, 1979, describes a process in which acetonitrile is prepared by high temperature reaction of carbon monoxide, hydrogen and ammonia over a transition metal catalyst. Various catalysts and operating conditions are exemplified.\nIt happens that the described reaction is usually accompanied by a number of side reactions, with the significance of particular reactions depending upon the catalyst and other conditions."}
{"text": "The scientific designing of systems for effecting controlled and sustained release of fragrances and malodour counteractants from semi-solid substrates including fabric, human skin and hair surfaces into the environment proximate such substrates has been the subject of ongoing investigations. Such systems include the targeted deposition onto substrate surfaces for extended periods of time of rupturable microcapsules composed of thin polymeric walls and containing substantial quantities of monophasic liquid compositions of such fragrance compositions and malodour counteractant compositions.\nIn designing and then employing such systems it has been found necessary to achieve optimization of (a) the nature and content of the targeted microcapsules, such as microcapsule dimensions, including average effective diameters and average wall thickness, polymer wall composition, specific ingredients, weight ratio of capsule wall: functional product composition, ratios of microcapsule volume:substrate area to which microcapsules are adhered and loaded, microcapsule weight:substrate area to which microcapsules are adhered in combination with (b) their respective rates of release of the functional products contained therein as a function of time, temperature and abrading use thereof, such as rubbing fabrics or skin having the microcapsules adhered thereto or, brushing of groups of hair follicles to which the microcapsules are adhered to.\nAccordingly, a need has been found to exist for (a) simulating such abrading use in the laboratory, and (b) enabling analytical measurement of the results of such abrading use on a continuing basis, including analytical measurement over an extended period of time of the several components of the headspace composition proximate the microcapsule-bearing substrate being abraded, as well as the rate of change of the several components of the headspace composition proximate the microcapsule-bearing substrate being abraded with respect to time and temperature.\nThe prior art discloses techniques for enabling the analysis of fragrances, including movement-activated fragrances adsorbed onto a porous substrate and evolving into the environment proximate such substrate, such as, towels as disclosed in U.S. Pat. Nos. 6,511,852 and 6,495,375. In addition, the use of relatively large and dense objects having a high degree of surface hardness, such as, steel spheres, for causing formation of dust from vinyl polymer-coated granules for the purpose of enabling analytical measurements, such as, analysis of vinyl polymer-coated enzyme granule dust is disclosed in U.S. Pat. Nos. 5,324,649 and 5,879,920 and published U.S. Patent Applications 2001/0056177 and 2002/0193275, wherein reference is made therein to the “Heubach attrition test”. The Heubach attrition test is described in detail in “Enzymes in Detergency”, ed. Jan H. Van Ee et al. (Marcel Dekker, N.Y., 1997) at pages 310-312 of Chapter 15 (Becker et al. “Formulation of Detergent Enzymes”). A schematic diagram of the Heubach apparatus is set forth on page 312, FIG. 3b. In addition, (i) for effecting particle size reduction, U.S. Patent Application 2001/0016259 A1, discloses the use of mobile steel balls maintained in motion by means of the operation of a shaker, and (ii) for effecting solubilization of solids in liquids, U.S. Pat. No. 5,672,456 discloses the use of a reciprocating shaker to hasten the dissolution process.\nHowever, nothing in the prior art discloses or infers any method or apparatus for effectively and efficiently enabling the qualitative and/or quantitative chemical analysis of the components of the headspace above abraded microencapsulates containing fragrances and/or malodour counteractants as a result of the simulated abrading use such as rubbing and/or brushing on such functional product-containing microcapsules which are affixed to a semi-solid substrate, including fabric, skin or hair follicles."}
{"text": "This application is based on and claims priority under 35 U.S.C. xc2xa7 119 with respect to Japanese Application No. 2002-285353 filed on Sep. 30, 2002, the entire content of which is incorporated herein by reference.\nThis invention generally relates to a walk-in apparatus for a vehicle seat. More particularly, the present invention pertains to the walk-in apparatus for the vehicle seat, which provides a forward slide movement of the vehicle seat relative to a vehicle floor in conjunction with a forward reclining movement of a seat back of the vehicle seat relative to a seat cushion.\nA known walk-in apparatus is disclosed, for example, in a U.S. Pat. No. 6,341,819 as a seat slide device for a vehicle.\nThe walk-in apparatus assembled to a slide rail mechanism, which slidably supports a vehicle seat in longitudinal direction relative to a vehicle floor, is adapted to release a lock plate by a released lever which rotates in conjunction with a pivotal movement of the seat back when the seat back of the vehicle seat is reclined as well as to release by a slide rail mechanism including an operation handle for unlocking the lock plate by hand. Such lock plate is applied biasing force for locking the slide movements of the slide rail mechanism.\nThus, such walk-in apparatus includes the reclining plate as a component of the reclining arm mechanism for rotating the seat back by a pivotal movement in forward direction of the vehicle, and the reclining plate are connected to release lever as a component of the slide rail mechanism through a rod.\nHowever, the known walk-in apparatus has the slide rail mechanism which is adapted to release the locking of the slide movement by way of both operating the operation handle and forwardly reclining the seat back. In this apparatus, if the seat back is reclined for operating the walk-in apparatus without knowing that the operating handle happens to be locked because of unexpected reasons, a member for locking the slide rail mechanism may be transformed or even broken due to input of the unreasonable load, as a result, the locking operation of the slide rail mechanism malfunctions.\nThus, the walk-in apparatus needs to be comprised in consideration of ensuring the locking operation of the slide rail mechanism even if the unreasonable load is applied to the member for locking the slide rail mechanism when the operating handle of the slide rail mechanism is locked.\nThe invention involves a walk-in apparatus of a vehicle seat which comprises an upper rail assembled to the vehicle seat, a lower rail engaged slidably with the upper rail and assembled to a vehicle floor; a lock plate locking a slide movement of the upper rail on the lower rail, a release lever which is adapted to engage with the lock plate for releasing the locking of the slide movement of the upper rail on the lower rail by the lock plate and a reclining plate which rotates in conjunction with a seat back when the seat back is rotated equal to or more than a predetermined angle in forward direction characterized in that the walk-in apparatus of the vehicle including a connecting member for connecting the release lever and the reclining plate is adapted to disengage the release lever from the reclining plate when a load is applied to the connecting member equal to or more than a predetermined value."}
{"text": "A classical computer operates by processing binary bits of information that change state according to the laws of classical physics. These information bits can be modified by using simple logic gates such as AND and OR gates. The binary bits are physically created by a high or a low energy level occurring at the output of the logic gate to represent either a logical one (e.g., high voltage) or a logical zero (e.g., low voltage). A classical algorithm, such as one that multiplies two integers, can be decomposed into a long string of these simple logic gates. Like a classical computer, a quantum computer also has bits and gates. Instead of using logical ones and zeroes, a quantum bit (“qubit”) uses quantum mechanics to occupy both possibilities simultaneously. This ability means that a quantum computer can solve certain problems with exponentially greater efficiency than that of a classical computer."}
{"text": "1. Field of the Invention\nThis invention generally relates to transporting timer protocols over a network using the Generic Frame Procedure. More specifically, the invention relates to an architecture and a procedure well suited to enable the transport of Sysplex Timer protocols over Generic Frame Procedure networks.\n2. Background Art\nOver the past several years and continuing to the present, computer manufacturers are providing processing architectures based on a multi-system shared data approach. Through these architectures, multiple large-scale computer systems, each of which is often referred to as a computer processing complex (CPC) or a central electronic complex (CEC), are interconnected, through, for example, a coupling facility or other inter-processor communication mechanism, to permit each such system to gain read-write access to data residing on one or more shared input/output devices, such as a direct access storage device (DASD). The resulting inter-connected computer system is commonly referred to as a “sysplex”. In a sysplex, as with a typical multi-processing environment, a processing workload is generally distributed among all of the inter-connected computer systems such that each computer system is responsible for processing a portion of the entire workload. Owing to the inherent high reliability and highly cost-efficient expansion potential of a sysplex architecture, sysplexes are particularly attractive in handling so-called critical business support applications that involve real-time transaction processing and can tolerate essentially no downtime.\nCertain currently available machines that can be readily incorporated into a sysplex, such as illustratively the Enterprise System/9000 (ES/9000) Series manufactured by the International Business Machines (IBM) Corporation, can each support, if appropriately configured, multiple actively and simultaneously executing copies of various operating systems (OS) to implement separate corresponding individual and unique application processing environments. (Enterprise System/9000 is a registered trademark, and ES/9000 is a trademark, of IBM Corporation)\nSysplexes, and other networks, may employ various technologies for transmitting data within the sysplex or network. For example, one technology, which has been in use for some time in telecommunication networks, but has only begun to find emerging applications in sysplex communications is referred to as SONET/SDH. SONET/SDH, which is an acronym for Synchronous Optical Network/Synchronous Digital Hierarchy, is a set of related standards for synchronous data transmission over fiber optic networks. SONET/SDH is currently used in wide area networks (WAN) and metropolitan area networks (MAN). A SONET system is comprised of switches, multiplexers, and repeaters, all connected by fiber. The connection between a source and destination is called a path.\nMore recently, the International Telecommunications Union (ITU) has proposed a new industry standard G.704.1 called Generic Frame Procedure (GFP). This is intended to allow standard datacom protocols such as ESCON or FICON to be encapsulated into a SONET/SDH compliant frame structure, so they can be transported across installed SONET networks. Since there is a large amount of SONET infrastructure installed by telecom carriers and other service providers, GFP is seen as the means for allowing enterprise systems to carry data traffic over existing SONET networks at very low incremental cost. In turn, this enables channel extensions over hundreds of kilometers for disaster recovery applications, such as large storage area networks. In order to remain competitive in the telecom carrier market, many wavelength division multiplexing (WDM) equipment manufacturers are adopting GFP transport.\nMany networks, however, are not fully compatible with GFP. For example, in a Sysplex, timing data, referred to as Sysplex Timer links, cannot be readily transported over GFP at the present time."}
{"text": "This invention relates to improvements made on the method of nitriding disclosed in the Japanese Pat. No. 776055 (Japanese Patent Publication No. 49-41023) of the same inventor. In this method, a salt-bath including titanium and/or zirconium as catalyst is used, and the electrolysis is made at relatively low temperature of less than 500.degree. C. Although nitrided case of considerably deep and high hardness is obtainable by this method, these characteristics are still confined to some extent, and besides that the method is only applicable to the low carbon, or low alloy-low carbon steels."}
{"text": "The present invention relates to a double layer capacitor, a method for the preparation thereof and the use thereof in a combined application with batteries.\nThe development of capacitor technology based on the principle of energy storage in the electrochemical double layer formed at the interface between an ion-conducting phase, i.e. the electrolyte, and an electron-conducting phase, i.e, the electrode, has provided capacitors of extremely high capacitance. Such capacitors are usually referred to as double layer capacitors, supercapacitors, electrochemical capacitor or ultracapacitors.\nThe energy storage in such capacitors may involve electrochemical processes like in batteries, i.e. energy is stored chemically in the capacitor electrodes and the electrode reactions involve redox processes. One such system is the rubidium silver iodide low voltage electrochemical Hypercap capacitor from Technautics. Alternatively, the capacitor may be based entirely on the double layer energy storage principle. Such capacitors will below be referred to as double layer capacitors.\nOne important application of double layer capacitors is in hybrid combinations with batteries. Traditionally, double layer capacitors provide high power capabilities compared to batteries, e.g. during pulse applications a higher power is obtainable from the double layer capacitor than from the battery. Accordingly, the power capability of hybrid combinations is higher than for the battery alone, and further the stress on the battery is reduced due to more uniform load. A number of patents describes the use of such battery-capacitor hybrid systems, among those U.S. Pat. No. 5,587,250 and U.S. Pat. No. 5,670,266 to Motorola.\nHigh capacitance and low impedance are the two main technical features providing the high power capabilities of double layer capacitors. The capacitance provides the energy for the load, whereas the low impedance allows good power accessibility in that the energy is available without too high internal losses.\nIn particular, as most pulsed loadings operate in the Hz-range from 1Hz to 1kHz, the capacitor impedance in this should be low, i.e. the main part of the energy stored in the capacitor should be available in this frequency range. A number of patents and patent applications describe approaches to double layer capacitors of low internal resistance.\nEuropean patent application EP 763,836 to Nisshinbo Industries discloses a polarisable electrode for use in a electric double layer capacitor, having a low internal resistance and a high capacitance, the carbon mixture of the electrode characterised by comprising fibrillated carbon. A conductive agent may be added to the carbon mixture to improve electrical conductivity. Such electrode has long life and can be charged and discharged at large electric currents.\nEuropean patent application EP 660,345 to Nisshinbo Industries discloses a polarisable electrode comprising a solid active carbon, characterised in that the electrode on the surface and/or inside has discontinuous portions free from said solid active carbon. A conductive agent such as graphite or carbon black may be added to the carbon mixture to improve electrical conductivity.\nU.S. Pat. No. 5,682,288 to Japan Gore-Tex, Inc. discloses a planar layered electrode comprising acicular electrically conductive particles to provide low electrical resistance and which has high electric storage capacity.\nU.S. Pat. No. 5,077,634 to Isuzu describes electrodes for double layer capacitors, which are compressed for reduction of internal resistance.\nLow impedance double layer capacitors should display low intra-component impedance as well as low inter-component impedance. Whereas the above prior art has been focussed on low intra-component impedance in the electrode structures, the impedance of the interface to the current collectors is just as important. Electrodes of high resistance at the current collector-electrode interphase, e.g. due to poor adhesion and poor electrical contact, will suffer from poor power capability due to internal losses.\nA number of patents describes double layer capacitors of elaborated electrode-current collector interface.\nJP-06,084,700 to Matsushita discloses a conductor layer made of carbonaceous material such as graphite provided between the electrode and the collector, the binder being an acrylic acid-styrene copolymer and providing increased adhesion for the electrodes. JP-61,102,025 to Matsushita discloses an electric double layer capacitor of polarising electrodes from activated carbon fiber felt, the electrode having a conductive layer on one side. The conductive layer may be aluminium, titanium, nickel or tantalum, stainless steel or a conductive paint containing carbon particles. The electrolyte of said capacitor is a solution of (NEt4) ClO4 in propylene carbonate. Upon repeated charging and discharging, the capacitor has stable capacitance.\nJP-60,050,914 to NEC discloses a double layer capacitor having a reduced internal resistance, the capacitor having electrodes covered by conductive layers consisting of a polymer such as polypropylene and a conductor such as carbon black or carbon fibers.\nU.S. Pat. No. 5,150,283 to Matsushita describes electrodes for a double layer capacitor composed of electrically conductive substrates, coated with a mixture of activated carbon with a water-soluble material-based binding agent. Such capacitor has low internal resistance, withstanding high voltages.\nU.S. Pat. No. 5,115,378 to Isuzu describes an electrode for a double layer capacitor formed from a porous sintered body of joined active carbon particles, conductivity being provided during sintering. Such capacitor has reduced internal resistance as low contact resistance between electrode body and current collector is obtained.\nU.S. Pat. No. 5,099,398 to Isuzu describes an electrode for a double layer capacitor applied on an electrically conductive film, the surface of which is dissolved in a solvent and the dissolved portions of said current collector being present in the pores of the electrode bodies. Such capacitor has low internal resistance due to intimate contact at the molecular level.\nAccordingly, several approaches to low internal resistance capacitors exist, including introducting of conductive coatings at the interface between the electrode and the current collector. Still, however, the is a need for a conductive coating which is stable under the severe conductions existing at this interface. This conductive coating should provide good adhesion to the current collector and to the electrode, and it should maintain its mechanical integrity upon handling such as folding and winding. Further, the coating should be stable at highly oxidative potentials at the positive electrode, and at highly reductive potentials at the negative electrode as well as being stable against corrosive electrolytes.\nTherefore, it is the objective of the present invention to provide a double layer capacitor with a conductive coating, which is mechanically, chemically, and eletrochemically stable during manufacture and operation. Such coating comprises a conductive agent, a binder and optionally auxiliary materials. Auxiliary materials may be materials used entirely during the processing of the conductive coating and which are removed from the final coating, or they be materials which facilitate the processing and which remain in the final coating.\nFrom a comprehensive study of a high number of conductive coatings it has now been found, that those coatings based on binders of the melamine resin type fulfil the above objectives.\nSurprisingly it was found, that the melamine resin binders provide long term reduction stability at potentials as low as +1.OV vs. Li/Li+. Whereas the considerable stability of melamine resins against oxidation is well known from their long term stability under ambient conditions, their high stability against reduction is surprising considering their chemical composition. Traditionally only fluorinated compound display stability under such reductive conditions.\nFrom the study it was also found, that despite the fact that melamine resins are hard and often brittle, the coating based on melamine resin binders showed good mechanical integrity and flexibility during processing of double layer capacitors.\nThus, according to the present invention a double layer capacitor is provided, which comprises a conductive coating comprising a melamine resin binder at the interfaces between current collectors and electrodes. In particular, the present invention provides a double layer capacitor comprising metal foil current collectors, carbon electrode structures with a polymer binder, conductive coatings comprising a melamine resin binder at the interfaces between current collectors and electrodes, and a non-aqueous electrolyte. Such double layer capacitors display long term low impedance at the electrode-current collector interface and high power capabilities.\nIn a preferred embodiment of the invention, the melamine resin is a alkylated melamine formaldehyde resin, preferably a methylated melamine formaldehyde resin. As reaction partner for the polymerisation process alkyd resins are preferred.\nFurther, the excellent performance has been found to be mostly pronounced for those double layer capacitors, which have conductive coatings of a thickness of 1-10 xcexcm.\nThe conductive coating comprising a melamine resin binder preferably has a composition prior to coating of:\n5-50% by weight, preferably 10-40% by weight, more preferably 20-35% by weight of carbon blacks;\n5-20% by weight, preferably 5-15% by weight, more preferably 10-15% by weight of graphite;\n5-40%, by weight, preferably 10-30% by weight, more preferably 15-25% by weight of melamine resin binder and polymerisation reaction partner;\n25-85% by weight, preferably 30-74% by weight, more preferably 35-53% by weight of solvent; and\n0-10% by weight, preferably 1-5% by weight, more preferably 2-5% by weight of auxiliary materials, preferably selected from the group consisting of dispersing agents, defoaming agents and rheological control agents.\nThe carbon blacks should display high structure and are advantageously selected from the group consisting of furnace blacks, acetylene blacks and lampblacks. The graphites should display low particle size, advantageously in the range of 0.5-20 xcexcm, preferably 0.5-10 xcexcm.\nThe solvents are preferably alcohols R1xe2x80x94OH, wherein R1 represents C1-C4 alkyl, and glycols and glycol ethers R2xe2x80x94(OCHR3CH2)nxe2x80x94OH, n=1-3, wherein R2 represents hydrogen or C1-C4 alkyl and wherein R3 represents hydrogen or methyl.\nThe dispersing agent may be non-ionic, anionic, cationic, as well as amphoteric. Preferably an anionic dispersing agent, such as Disperbyk 170 from BYK Chemie is used.\nThe defoaming agent may be a mineral oil or silicone oil defoaming agent, preferably a silicone oil defoaming agent, such as BYK-080 from BYK Chemie.\nThe rheological control agent may be an organo clay, silica and castor oil derivative, preferably an organo clays, such as Viscogel B7 from Chimica Mineraria SpA.\nThe conductive coating may be commercially available products with a melamine resin binder, such as the XZ302 screen printing dye from Wiederhold Siebdruckfarben of Germany. This dye is based on carbon black and graphite and a thermosetting resin of the melamine resin type, the solvent mixture comprising 2-(2-butoxyethoxy) ethanol, 2-butoxyethanol, butan-1-ol and 1-methoxypropan-2-ol. The retarder UV4 from Wiederhold Siebdruckfarben may be used for optimisation of viscosity and drying time.\nIn a preferred embodiment of the invention the melamine resin-based conductive coating is used in double layer capacitors of electrolytes based on tetraalkyl-ammonium salts. Thus, it has been found, that those double layer capacitors based on tetraalkylammonium salts have a high capacitance and a higher power capability than double layer capacitors using electrolyte compositions of other salts like lithium and sodium salts. Although not fully understood, the formation of any interface layer appears highly dependent on the ionic species of the electrolyte, as well as on the surface structure and chemical composition of the electrode carbon. The excellent performance of double layer capacitors based on tetraalkylammonium salts is ascribed to little, thin, stable and dense interface layers being formed at the electrode-electrolyte interfaces of such capacitors, allowing a narrow charge separation and a high capacitance. In contrast, in the case of lithium salts, thicker, unstable and less dense interface layers are formed, which provide less capacitance. Accordingly, the concept of conductive coatings based of melamine resin binder is particularly advantageous in those cases, where the electrolyte salt is a tetraalkylammonium salt. Thus, the higher capacitance of capacitors of such salts is most effectively applied when combined when the low impedance and low internal loss of the conductive coatings of the present invention.\nAlthough a high number of tetraalkylammonium salt may be used, the above advantage is especially pronounced for tetramethylammonium tetrafluoroborate, tetraethylammonium tetrafluoroborate, tetrabutylammonium tetrafluoroborate, tetramethylammoniumhexafluorophosphate, tetraethylammonium hexafluorophosphate and tetrabutylammonium hexafluorophosphate, in particular tetraethylammonium tetrafluoroborate. The electrolyte solvent may be any non-aqeuous solvent selected from the groups of carbonates, lactones and nitriles. Preferably the electrolyte solvent is selected from the group of ethylene carbonate, propylene carbonate, dimethyl carbonate, diethyl carbonate, ethyl methyl carbonate, xcex3-butyrolactone, xcex3-valerolactone and acetonitrile and mixtures thereof. More preferably the electrolyte solvent is xcex3-butyrolactone.\nIn a preferred embodiment of the invention the electrolyte is confined in a separator, consisting of a porous structure made from a polymer, preferably polyethylene or polypropylene.\nIn a preferred embodiment of the invention the current collector is a metal current collector, preferably of nickel, copper or aluminium, more preferably of aluminium.\nThe electrode structures of the double layer capacitors of the present invention comprises carbon, binder, solvent and optionally graphites and any auxiliary compounds\nIn a preferred embodiment of the invention the carbons of the electrode structure are high surface active carbons.\nIn a preferred embodiment of the invention the binder of the electrode structure is selected from the group consisting of melamine resins, polyvinyl butyrals and fluorocontaining polymers, preferably PVdF and PVdF-copolymers, including PVdF-hexafluoro propylene copolymers. For the processing of such binders, solvents therefor are preferably selected from the group consisting of glycol ethers and glycol ether esters, dimethyl formamid, dimethyl acetamid and N-methyl-pyrrolidone.\nThe invention further relates to a simple, low-cost method of preparing double layer capacitors.\nAccording to the invention the conductive coating is prepared from a paste comprising carbon black, graphite, melamine resin binder, solvent and optionally any auxiliary compounds. Carbon blacks and graphites materials are added along with any dispersing agents and defoaming agent and a polymerisation reaction partner are mixed in alcohols, glycols or glycol ethers. The melamine resin is added along with any rheological control agent to produce a uniform paste, which is coated or printed onto the current collector and heated to 100-150xc2x0 C. for 10-30 min.\nElectrode pastes are prepared similar to the above conductive paste from active carbon, binder, solvent and optionally graphite and any auxiliary compounds, and coated or printed onto the conductive coating.\nThe coating or printing technique applied is preferably a screen printing, gravure printing or a slot die coating technique.\nDouble layer capacitors are formed by sandwiching between two of the above conductive coating-electrode laminates a porous separator. The capacitor laminate is subsequently confined in a polymer coated metal pouch and impregnated with the electrolyte solution. Eventually, the pouch is sealed.\nA final aspect of the invention relates to the use of double layer capacitors according to the present invention in a hybrid combination with batteries, in which combination the double layer capacitor provides peak current, whereas the battery provides the base load currents. In such applications the battery capacity utilisation is improved, as the stress on the battery is reduced."}
{"text": "The power FET has become a limiting component in many power management circuits used to regulate voltage and/or current in DC-DC converters, AC-DC inverters, or AC transformers. In recent decades, significant efforts have been made to boost the power FET's output currents and switching speeds to keep pace with the integrated circuit “transistor shrinks”. The smaller transistors increase system switching speeds and the transistor densities within an integrated circuit. Higher densities allow vastly more transistors to be incorporated into a single semiconductor chip, causing it to require larger operational currents. Similarly, smaller transistor dimensions also increase operational system speeds. Any inability to supply a combination of larger currents at higher switching speeds leaves many systems “under powered”. This is particularly the case in multi-core microprocessors where the inability to supply required current at suitably higher speed impairs reliable data transfer with external memory circuits. As a result of this deficiency, processor cores will typically operate at 25%-30% utilization rates. This problem becomes more acute when processor cores are configured in parallel. For instance, a 16 core microprocessor array will function slower than a 4 core microprocessor array due to in adequate power refresh cycling. Therefore, it is desirable to increase power efficiencies by developing power FETs that enable arbitrarily high currents to be switched at arbitrarily high switching speeds. While higher switching speed and current levels are of specific benefit to power FETs, methods that enable switching speeds above 400 MHz or over a distinct band of frequencies could be usefully applied to many other FET applications. Therefore, generalized methods to tune a FET's switching speed and/or current output is also desirable.\nCo-location of high efficiency switched-mode power management devices with one or more processor cells also reduces the overall system power losses through much shorter interconnect circuitry. Methods and apparatus that improve the efficiency of supplied power to a processor core by co-locating power management in immediate proximity to computational die are therefore enabling and desirable to the enhanced utilization of microprocessor arrays and the improved operational efficiency of high-speed computational systems.\nMost approaches to increase power FET switching speeds have focused on reducing the transistor gate capacitance. This is done generally by making the gate electrode smaller, and by using more sophisticated electronic doping configurations within the transistor junction. A significant drawback to these elements of the prior art is the generation of larger quantities of heat, which is undesirable and must be properly managed to safeguard the performance of neighboring semiconductor devices. While smaller gate structures lower capacitance, they also increase current densities flowing through the transistor junction to attain high currents, which, in-turn, increases generated heat to levels that prevent co-location of power management in close proximity to computational semiconductor die. Higher current levels often require thermal management mechanisms to be added to the system to drain the excess heat generated by the higher levels of resistive loss. Additional or more sophisticated thermal management requirements increase the cost and design complexity of the overall system. Prior art solutions advance improved doping topologies to reduce the ON-resistance of the power FET's transistor junction. However, the proposed solutions alone do not reduce ON-resistance to sufficiently low levels that allow power management systems to be monolithically integrated with or placed in close proximity to the active system devices, or to mitigate or eliminate thermal management devices from the power management system. This is a particular problem in power management systems that are assembled from discrete components. Therefore, it is also desirable to reduce the ON-resistance of a power FET by orders of magnitude to levels that mitigate or eliminate thermal management systems in a power management device.\nMany power management systems are utilized in mobile platforms that are subject to frequent or unpredictable mechanical shock. Solder joints are used to electrically interconnect surface mounted passive components (resistors, capacitors, inductors) on a printed circuit board in electrical communication with one or more semiconductor die. Lead-free solders used to achieve modern environmental standards do not have the mechanical integrity of the lead-based solders they are replacing. Fractures in solder joints are the dominant cause for field failures in mobile systems. Additionally, solder joint failures in power management devices are the leading cause for grounding aircraft due to unscheduled maintenance. Therefore, methods that eliminate solder joints from a power management system by monolithically integrating all components on to a semiconductor die are also beneficial and a desirable objective of the present invention."}
{"text": "A dual-mode mobile terminal is a mobile terminal capable of supporting two radio technologies having different communication schemes, and is mainly used in a region in which heterogeneous communication networks coexist. As a representative example of the dual-mode mobile terminal, many developers and companies are conducting intensive research into mobile terminals configured to use not only Long Term Evolution (LTE)-based wireless communication but also Code Divisional Multiple Access (CDMA)—based wireless communication. An enhanced High-Rate Packet Data (eHRPD) scheme may be used as a representative example of CDMA. The eHRPD is a new version of the 1xEV-DO upper layer protocol stack developed by the 3GPP2 standards committee so as to enable wireless communication network interoperability for LTE.\nAlthough a dual-mode terminal capable of communicating with the LTE network and the eHRPD network according to the present invention will hereinafter be disclosed as an example for convenience of description and better understanding of the present invention, it is obvious to those skilled in the art that other wireless communication schemes can also be applied to the present invention without difficulty.\nIf a user equipment (UE) supports two different communication schemes, the UE may prefer to receive services through a high radio technology from the viewpoint of a general communication enterprise. That is, assuming that the UE supports the LTE scheme and the eHRPD scheme, the UE may prefer to receive a necessary service through the LTE network. However, service coverage of the eHRPD network is larger than that of the LTE network at an initial service step in which switching from the legacy eHRPD network to the LTE network begins. In the initial service step in which switching from the legacy eHRPD network to the LTE network begins, the service coverage of the eHRPD network is larger than that of the LTE network. If the UE is located at a coverage boundary of the LTE network under the above-mentioned situation, the UE may be handed over to the eHRPD network by a network configured handover condition. On the contrary, if the UE receives a necessary service through the eHRPD network, the UE periodically searches for an LTE cell. Then, the UE having searched for the LTE cell may attempt to immediately access the LTE cell without network control. In this case, if the UE is located at an area with poor coverage, such as a cell edge of the LTE cell or a shaded region, the UE again determines the network configured handover (HO) condition. If the network configured handover (HO) condition is satisfied, the UE is handed over to the eHRPD network. The above-mentioned phenomenon (i.e., ping pong phenomenon) may repeatedly occur in so far as the UE stays in a specific region in which the corresponding condition is satisfied.\nThe above-mentioned problem is based on the fact that a first condition in which the UE performs handover from the LTE network to the eHRPD network is different from a second condition in which the UE performs handover from the eHRPD network to the LTE network. Accordingly, the present invention defines a handover (HO) execution condition capable of solving the above problem, and proposes a method and apparatus for performing efficient handover on the basis of the defined HO execution condition."}
{"text": "Fuel cells have been used as a power source in many applications. For example, fuel cells have been proposed for use in electrical vehicular power plants to replace internal combustion engines. In proton exchange membrane (PEM) type fuel cells, hydrogen is supplied to the anode of the fuel cell, and oxygen is supplied to the cathode. PEM fuel cells include a membrane electrode assembly (MEA) comprising a thin, proton transmissive, non-electrically conductive solid polymer electrolyte membrane having the anode catalyst on one of its faces and the cathode catalyst on the opposite face. The MEA is sandwiched between a pair of electrically conductive elements which serve as current collectors for the anode and cathode, and contain appropriate channels and/or openings therein for distributing the fuel cells gaseous reactants over the surfaces of the respective anode and cathode catalyst. A typical PEM fuel cell and its membrane electrode assembly (MEA) are described in U.S. Pat. Nos. 5,272,017 and 5,316,871 issued respectively Dec. 21, 1993 and May 31, 1994 and assigned to General Motors Corporation.\nThe term “fuel cell” is typically used to refer to either a single cell or a plurality of cells depending on the context. A plurality of individual cells are commonly bundled together to form a fuel cell stack. Each cell within the stack comprises the MEA described earlier, and each such MEA provides its increment of voltage. Typical arrangements of multiple cells in a stack are described in U.S. Pat. No. 5,763,113 assigned to General Motors Corporation, and which is herein incorporated by reference.\nDuring fuel cell system start-up, a compressor is often utilized to provide compressed air or oxygen to the fuel cell cathode inlet while hydrogen is introduced to the fuel cell anode inlet. Many fuel cell systems currently use large battery packs to start and operate the compressor before output power is available from the fuel cell stack. The battery pack often is utilized with various DC/DC converters to provide the high voltage necessary for the compressor. Usually, a DC/DC converter is needed to get the battery voltage up to the stack level, then another DC/DC converter is needed to get the voltage from the stack level up to the compressor motor controller level. The large battery pack and DC/DC converters contribute significantly to the weight, volume, and cost of the fuel cell system. Thus, it is desirable to provide a simplification of the fuel cell system as well as reducing the mass volume and cost of the system."}
{"text": "The term “packaging” applied to diode-laser bars refers to mounting a diode-laser bar on some sort of cooling base or heat-sink. This base may be a relatively massive base providing a “conductively cooled package” (CCP). For higher power operation, the base may be water-cooled, for example by a micro-channel arrangement. Typically the diode-laser bar is soldered “p-side down” onto an insulating sub-mount. The sub-mount, in turn, is soldered to the heat-sink.\nPackaging involves a number of compromises and trade-offs. As the materials of the bar, sub-mount, heat-sink, and solder have different thermal properties, particularly the coefficient-of-thermal-expansion (CTE), there is invariably a trade-off between thermal performance and CTE-matching (or mismatching). Indium (In) soldered bars on copper (Cu) perform well thermally, i.e., tolerate some CTE mismatch, but soft solders such as the indium solder can be unreliable. One hard solder that is essentially completely reliable is an 80:20 gold-tin (Au—Sn) solder. This solder however cannot tolerate any CTE mismatch between the soldered components.\nThere is a need for a diode-laser bar packaging arrangement that will provide effective CTE-matching using high thermal-conductivity materials, Preferably the packaging arrangement should be capable of being water-cooled."}
{"text": "Technical computing applications and/or programming environments provide operators, (e.g., engineers, scientists, software developers, mathematicians, educators, etc.), across a diverse range of industries, with an environment for performing technical computing. The technical computing applications and/or programming environments provide mathematical and graphical tools for mathematical computation, data analysis, visualization and algorithm development. The technical computing applications and/or programming environments may be used to solve complex problems by developing mathematical models that can be used to simulate and/or analyze the problems.\nThe technical computing applications and/or programming environments are often implemented as desktop applications that allow operators to interactively perform complex analysis and modeling on a single workstation. However, the single workstation can be limiting to the size of the problem that can be solved due to limited processing capacity associated with the single workstation or an amount of time available to solve the problem. In this case, the single workstation may be allocated full-time to performing the analysis and modeling while other computing resources, on a network, remain idle.\nHence, distributed computing is useful for distributing tasks to multiple computing resources, especially for a computationally intensive problem. However, most technical computing applications or programming environments were not initially designed for distributed computing. Additionally, developing new technical computing applications for use in a distributed computing environment may be prohibitively costly, time consuming, and/or complex. Thus, users may resort to performing technical computing using a desktop application, on a single workstation, which may limit a level of complexity of problems to be solved and/or increase processing time relative to performing the technical computing in a distributed computing environment."}
{"text": "The term \"branding\" means the marking of timber with identification symbols relating to the type or quality (grade) of the timber. After the timber has been felled it is normal to rough saw the green timber before it is dried. The green timber has a high moisture content (e.g. greater than 100% for softwoods and up to 150% for hardwoods) and must be dried before use, usually outdoors or in a kiln. Green timber is particularly difficult to brand on account of the high water content. Moreover hardwood is generally wetted during sawing as the saw blade often has to be cooled with a water spray.\nExisting methods of branding freshly-sawn timber involve the use of various types of ink which provide neither distinctly legible grades of mark nor marks which can be maintained on the timber for a sufficient length of time, particularly in view of the rigorous conditions under which sawn timber is transported, stored and worked.\nIn particular, the branding must show satisfactory water resistance for storage outdoors. As timber is quickly passed out of the saw mill the brand must also be capable of drying within a short time. Typically less than 5 to 10 minutes. The branding fluid must also be economical and compatible with existing equipment.\nIt is an object of the present invention to mitigate these disadvantages associated with known methods of branding timber."}
{"text": "1. Field of the Invention\nThe present invention relates to techniques used in the field of microbial industry. More precisely, the present invention relates to a method for producing L-lysine or L-arginine by fermentation and a microorganism used in the production method.\n2. Description of the Related Art\nAmino acids such as L-lysine, L-glutamic acid, L-threonine, L-leucine, L-isoleucine, L-valine and L-phenylalanine are industrially produced by fermentation using microorganisms that belong to the genus Brevibacterium, Corynebacterium, Bacillus, Escherichia, Streptomyces, Pseudomonas, Arthrobacter, Serratia, Penicillium, Candida or the like. In order to improve the productivity of these microorganisms, strains isolated from nature or artificial mutants thereof have been used. Moreover, various techniques have been disclosed for increasing the L-amino acid producing abilities by using recombinant DNA techniques to enhance L-amino acid biosynthetic enzymes.\nProductivities of L-amino acids have been considerably increased by breeding of microorganisms such as those mentioned above and improvements of production methods. However, in order to respond to further increase in demands in future, development of methods for more efficiently producing L-amino acids at lower cost have still been desired.\nAs methods for producing L-amino acids by fermentation of methanol, which is a fermentation raw material available in a large amount at a low cost, there have hitherto known methods using microorganisms that belong to the genus Achromobacter or Pseudomonas (Japanese Patent Laid-open (Kokai) No. 45-25273), Protaminobacter (Japanese Patent Publication (Kokoku) No. 49-125590), Protaminobacter or Methanomonas (Japanese Patent Laid-open No. 50-25790), Microcyclus (Japanese Patent Laid-open No. 52-18886), Methylobacillus (Japanese Patent Laid-open No. 4-91793), Bacillus (Japanese Patent Laid-open No. 3-505284) and so forth. The inventors of the present invention have so far developed methods for producing L-amino acids using Methylophilus bacteria based on breeding by artificial mutagenesis and recombinant DNA techniques (Japanese Patent Application No. 11-368097).\nIn recent years, there have been identified proteins that have a function of specifically excreting an L-amino acid to outside a cell of microorganism and genes therefor, and in particular, Vrljic et al. identified a gene involved in excretion of L-lysine from a Corynebacterium bacterium to outside of a cell (Vrljic M., Sahm H., Eggeling L., Molecular Microbiology 22:815–826 (1996)). This gene was designated as lysE, and it was reported that L-lysine producing ability of Corynebacterium bacteria could be improved by enhancing this gene in Corynebacterium bacteria (WO97/23597). It is also known that productivities for some L-amino acids can be improved by increasing expression amounts of amino acid excreting proteins in Escherichia coli (Japanese Patent Laid-open No. 2000-189180). For example, it is reported that productivities of cysteine, cysteine and so forth can be improved by enhancing expression of ORF306 gene in Escherichia coli (EP885962).\nHowever, there has so far been disclosed no example of demonstrating that the excretion process of an amino acid constitutes a serious obstacle for amino acid production by fermentation of methanol using a methanol assimilating bacterium. There is also no report on any amino acid excretion gene that can provide such an excretion activity in a methanol assimilating bacterium.\nFurthermore, it has not been known that the lysE gene has a function of excreting amino acids other than L-lysine."}
{"text": "1. Technical Field\nThe present disclosure relates to a component and method for producing a computerized display output wherein information bearing hexagons having common borders are rendered so as to appear above or below an underlying map.\n2. Description of the Related Art\nIn the field of computer graphics, image data is output to a display device to provide information to the computer user. The generation of images or objects to be represented on the display device is conventionally referred to as rendering.\nIn some computer images, the information displayed is map data, which is useful to the computer user in many circumstances. For example, in some cases, the map data identifies geographic boundaries, natural or man-made structures, topography, and the like. In other cases, the information on the map is combined with additional information such as motion and weather. Overall, maps rendered on computer displays are very useful to the general public."}
{"text": "The present invention relates to a bulky asphalt-impregnated sheet having different properties on both surfaces thereof and a method and a system for manufacturing such bulky asphalt-impregnated sheet.\nRecently, we have developed an asphalt roofing which comprises a base material of a bulky sheet of non-woven fabric made of filaments intertwisted with each other and an asphalt with which said base material is uniformly and completely saturated and which also covers both sides of said base material. This asphalt roofing can be produced by feeding said bulky sheet along a vertical path into a bath of molten asphalt, thereby permitting the air contained within the bulky sheet to successively escape vertically through the bulky sheet itself into the atmosphere while permitting the bulky sheet to be uniformly and completely saturated with said molten asphalt, and then withdrawing said sheet into the atmosphere and drying the same. The present invention utilizes an intermediate product of such asphalt roofing, namely, a bulky sheet which is impregnated with and covered, on its both sides, by asphalt which is still in a fluid or semi-solidified state, to produce a novel product, that is a bulky asphalt-impregnated sheet having materials of different properties and different functions on both sides thereof. This novel product can be produced by dipping a bulky sheet of non-woven fabric into a bath of molten asphalt to saturate said bulky sheet with the asphalt, withdrawing said sheet upward from said bath of molten asphalt, scraping off the asphalt on one surface, either front or rear surface, of said sheet while the sheet is being pulled upward and the asphalt is held in a fluidized or semi-solidified state, and then supplying and firmly adhering a material of different property onto said sheet.\nThe conventional asphalt roofing is produced by preparing a felt-like base material mixed with paper or asbesto fibers or a base material made of non-woven fabric of synthetic fibers, impregnating said base material with molten asphalt and then forming asphalt coating layers on both sides of said base material. In case of the felt-like base material made of paper or asbesto fibers, the base material has a thickness of about 1.4 mm, including base material having a thickness of about 0.6 mm and a surface asphalt layer having a thickness of about 0.4 mm. Therefore, the scraping off of the thin surface layer of the asphalt under fluidized or semi-solidified state will necessarily result in breaking or tearing of the base material made of the fragile fibers.\nAlso, in the case of a base material made of non-woven fabric, the base material has a thickness below 2 mm and the surface asphalt layer has a thickness to below 0.5 mm. The base material has a sufficiently high strength to allow the surface asphalt layer to be scraped off, but the exposed surface is hard and smooth so that the adhering ability is decreased. Furthermore, since the base material has a small thickness, it tends to soften and deform if the material of different property is supplied at high temperature. Accordingly, a roofing having materials of different properties on both side surfaces cannot be obtained, except a combined structure including two sheets adhered together."}
{"text": "Coupling devices have been used extensively for many years in the railroad industry to connect adjacent ends of a pair of railway cars together to form a train consist of several cars. On those railway cars which will be used in interchange service, however, such coupling devices must receive approval from the Association of American Railroads (AAR) before they can be installed on the cars. It has come to be quite well known, in this railway application, that such coupling devices perform a number of important functions. Obviously, an important function of the standard coupling device is that it facilitates the connection and disconnection of individual cars to and from, respectively, a train consist. Another important function of a standard coupler is that it enables such railway cars to negotiate the various curved portions of the track structure which are encountered during operation. Further, these coupling devices will enable a number of cars to be readily combined so that a train consist can be made up, or such cars can be easily separated individually for either loading or unloading cargo thereto or therefrom, respectively, as required. This will allow a railroad to leave a particular car at a customer's plant while delivering or picking up other cars at other locations.\nIn more recent times, the railroad industry has come to recognize, however, that a number of rather significant advantages can be achieved by the interconnection of several railway cars together to form a generally semi-permanent unit. This is particularly the case, for example, where such individual railway cars are adapted for use in what is commonly known in the railroad industry as \"piggyback\" service. A primary reason for this is that the cargo to be either loaded or unloaded is brought to or removed from, respectively, a central location which is owned and operated by the railroad normally. Generally, this cargo is either over-the-road trailers or very large containers to be shipped by sea. The individual cars which have been connected together in such substantially semi-permanent fashion are commonly known in the railroad industry as a \"10-pack\". These 10-pack units do not require the use of standard couplers except at each outer end of each 10-pack unit. The primary reason why such standard coupling devices are not required is that in view of their dedicated service these 10-pack units will only be broken periodically. For example, this will normally only occur when maintenance must be carried out on an individual coupler component or on some other component positioned on such railway cars that will necessitate such car be taken out of service at least temporarily. It has become obvious to the railroad industry that with the use of this coupling arrangement it is possible for them to achieve a significant reduction in their costs. Such cost-savings can be attributed to a number of reasons. For example, these reasons include lower equipment weight which results in enhanced energy savings and fewer railway trucks which results in both lower equipment cost as well as a significant reduction in maintenance requirements. However, now with the rather extensive use of these semi-permanent coupling arrangements, particularly with new cars being built for piggyback service, and with the ever increasing loads that are now being carried by modern railway cars and trains, it has been determined that it is of the utmost importance for a close-buttoned relationship to be maintained between the draft components of the coupler devices. Such close-buttoned relationship is necessary, for example, so that the effects of the impact forces, which are usually encountered under normal buff conditions during train operation, can be reduced to an acceptable level. In this manner, unnecessary damage to both cargo and rolling equipment can be held to a minimum.\nTaught in U.S. Pat. No. 4,258,628, is one prior art type of articulated coupling apparatus which can be used for the above-identified purpose of connecting adjacent ends of a pair of railway cars together in such semi-permanent manner. This particular articulated coupling arrangement, as taught and illustrated therein, includes a male connection member secured to one end of a first railway car body and a female connection member which is secured to an adjacent end of a second railway car body. However, as can be seen the ends of both the male connection member and the female connection member which are secured to such car bodies do not provide any flexibility in how they are attached to such car body. The female connection member is rotatably-engaged in a center plate bowl portion of the bolster of a railway car truck in this prior art arrangement. This rotatable engagement is carried out in a manner that is well known in the railway art. The outer end portion of the male connection member is disposed for movement within a cavity formed in the outer end portion of such female connection member. To connect both the male connection member and the female connection member together in such semi-permanent manner, a pin member is utilized. This pin member is positioned in a vertical direction and is disposed in axially aligned apertures which are formed in each of the male connection member and the female connection member.\nAs taught in this reference, the aperture which is formed in such male connection member for receiving the pin member therein must be somewhat larger than the pin member itself. This larger aperture is required so that, while in operation, certain required movements of the coupling arrangement can be achieved. Additionally, a rear surface portion of such aperture that is located in such male connection member and which receives the pin member therein has a horizontally disposed concave configuration and a vertically disposed convex configuration. This particular configuration is required because it enables both the male connection member and the female connection member to move in each of a horizontal direction and a vertical direction in relationship to one another. Furthermore, this configuration provides, at the same time, a relatively substantial area of surface contact between the rear surface area of such pin aperture and the pin member itself.\nThe outer end surface of the outer end portion of such male connection member is provided with a convex configuration. This convex configuration abuttingly engages a complimentary concave surface which is formed on a front face portion of a follower member. In this coupling arrangement, such follower member is carried within the rear portion of such cavity disposed in the outer end portion of such female connection member. A pair of vertically disposed, slot-like cavities are provided on such follower members adjacent the rear surface thereof.\nA first portion of a resilient member is disposed within each of such vertical slot-like cavities. A second portion of each such resilient member extends outwardly from such rear face portion of the follower member. In this manner, a portion of the exposed outer surface of each such resilient member can be engaged by a vertically disposed wedge-like element. This wedge-like element must be provided in this prior art coupling arrangement to urge both the follower member and the male connection member forward.\nWhen this occurs the rear surface portion of the aperture formed in the outer end of the male connection member will be maintained substantially in contact with the pin member at all times.\nThis contact between such pin member and the rear surface portion of such aperture in the male connection member is absolutely necessary in this coupling arrangement because most of the articulated connecting members are manufactured as cast components. Consequently, in order for the manufacturer to achieve any significant reduction in the cost of this coupling device, such cast articulated connecting members receive very little, if any, finish-type machining which could provide the necessary and desirable dimensional control. In other words, these cast connecting members will generally be assembled into such coupling device as cast. Therefore, as a result of this cost-saving practice, it is quite often very difficult to provide an articulated coupling device which will be self-adjusting under the various wear conditions that will be encountered by such coupling device during in-track operation. Nevertheless, it is of the utmost importance to minimize the slack encountered in the various coupling connections during such in-track service for the reasons discussed above.\nOther prior art-type articulated coupling arrangements are taught in U.S. Pat. No. 3,716,146 and Canadian Patent number 1,231,078. However, the shortcomings which are overcome by the invention to be taught hereinafter are also found in these prior art references."}
{"text": "1. Technical Field\nThe present disclosure relates to a surgical fastener applying apparatus and, more particularly, to a surgical fastener applying apparatus having reusable and disposable components.\n2. Discussion of Related Art\nSurgical fastener applying apparatus, wherein tissue is first grasped or clamped between opposing jaw structures and then joined by means of surgical fasteners, are well known in the art. In some such apparatus, a knife is provided to cut the tissue which has been joined by the fasteners. The fasteners are typically in the form of surgical staples, although, other surgical fasteners may also be utilized, such as, for example, clips or two part polymeric surgical fasteners.\nSurgical fastener applying apparatus typically include two elongated beam members which are used to capture or clamp tissue therebetween. Typically, one of the beam members carries a disposable cartridge assembly which houses a plurality of staples arranged in at least two lateral rows, while the other beam member includes an anvil which defines a surface for forming the staple legs as the staples are driven from the cartridge assembly. Where two part fasteners are used, the beam member which includes the anvil carries a mating part of the two part fastener, e.g. the receiver. Generally, the staple formation process is affected by the interaction between one or more longitudinally moving camming members and a series of individual staple pushers. As the camming members travel longitudinally through the cartridge carrying beam member, the individual pusher members are biased upwardly into a backspan of the staples supported within the cartridge assembly to sequentially eject the staples from the cartridge. A knife may be provided to travel with the camming members between the staple rows to cut the tissue between the rows of formed staples. An example of such an instrument is disclosed in U.S. Pat. No. 7,631,794, which is incorporated herein in its entirety by reference.\nBecause of the dangers associated with improper sterilization of surgical apparatus, fastener applying apparatus are typically disposable after use. Although the cartridge assembly may be replaced to perform multiple fastener applying operations on a single patient, the staple applying apparatus is typically disposable after a surgical procedure has been completed. This requirement of disposability may increase the costs associated with surgical procedures. Although reusable fastener applying apparatus have been developed, such apparatus can be overly complex and prove difficult to sterilize.\nA need exists in the art for a fastener applying apparatus which includes reusable components, is not overly complex and is configured to facilitate proper sterilization after use in a surgical procedure."}
{"text": "In a receiver of a wireless communication system, symbols in a single code block transmitted to a transmitter during a single Transmit Time Interval (TTI) experience different channels due to frequency selective fading of the channels. Accordingly, Post Processing Signal to Interference-plus-Noise Ratios (PPSINRs) of the symbols in the single code block have different values. An effective SINR indicating a reception quality for the single code block may be obtained by combining the PPSINRs of all of the symbols in the code block.\nIn general, methods for calculating an SINR of each symbol in a code block in a Single Input Single Output (SISO) system are explicitly defined. However, in a Multiple Input Multiple Output (MIMO) system, it is not as easy to calculate an SINR for each symbol because a signal received through a specific antenna contains an interference signal generated by a different antenna. In particular, since a Maximum Likelihood (ML) detector, which is used to improve performance at a MIMO receiver in currently, joint-detects the interference signal and thus does not treat the interference signal as a simple noise. Therefore, it is not easy to obtain an SINR for a symbol of a specific stream when the ML detector is used.\nAccordingly, various related-art methods have been provided to obtain a link quality for each stream in an approximate method without calculating an SINR when a ML receiver for each symbol of a stream is used. For example, a method for measuring a link quality in a Minimum Mean Square Error (MMSE) detection method in a receiver using the ML detector has been provided. However, this method has a problem that the link quality is not exactly measured when a Signal to Noise Ratio (SNR) is high. That is, when the SNR is high, the PPSINR of each symbol in the code block is theoretically very low because power of the interference signal generated by the different antenna is operated as a noise. However, since the ML detector actually joint-detects the interference signal, the PPSINR of each symbol in the code block is high. Therefore, an error occurs in the result of measurement of the link quality when the MMSE detection method is used. Such a link quality measurement error may result in an error in channel status information feedback such as Channel Quality Indicator (CQI) feedback, Rank Indication (RI), or Precoding Matrix Indicator (PMI), and thus may cause grave deterioration in performance of a link throughput.\nTherefore, there is a demand for a method for measuring a link quality without causing a measurement error or performance deterioration when an ML detector is used in a MIMO system\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "1. Field of Invention\nThe present invention relates to an ink ribbon cartridge containing a replaceable wide ink ribbon.\n2. Description of Related Art\nGenerally an ink ribbon cartridge is used because of its simplicity of handling when a thermal printer is used to print on common printing paper. However, when a thermal printer (e.g., a line printer) uses a large ink ribbon cartridge, it becomes uneconomical if the whole cartridge thrown away after use. To obviate this drawback, the ink ribbon alone is replaced.\nFor this type of ink ribbon, there is employed an ink ribbon including a band-like sheet comprising a resin film with a layer of ink formed on one side of the ribbon and wound around a pair of paper spools. A prior art ink ribbon cartridge using this type of ink ribbon was constructed such that, at the time of ink ribbon replacement, the cover is opened, the used ink ribbon is taken out, and four spindles inserted in both ends of a pair of paper spools of the ink ribbon are removed and then placed in both ends of one pair of paper spools of a new ink ribbon. After the insertion of the new ink ribbon into the cover, the cover is closed. Subsequently, the ink ribbon cartridge is mounted in a printer or a facsimile system as supported by a spindle on bearings that transmit the torque to the paper spools through the spindles, thereby winding the sheet around the paper spool.\nHowever, the above-described prior art ink ribbon cartridge is of such a structure that any one of the four spindles can be inserted into either of both ends of a pair of paper spools. In addition, in any combination of spindles and spools, the ink ribbon is inserted into the cover to allow the cover to close. Therefore, failures in printing due to improper mounting of the ink ribbon often occurred.\nThat is, since a layer of ink is formed on one side of the sheet, the ink cannot be transferred to a recording paper if the sheet is mounted with the wrong side out. Also, if the direction of winding is reversed, the sheet can not be wound up, resulting in a failure in printing."}
{"text": "1. Technical Field\nThis invention relates in general to semiconductor memory devices and, more specifically, to a cache tag that can be configured in accordance with a selected burst length.\n2. State of the Art\nModern memory systems for personal computers and the like generally include a main memory that consists of approximately 32 Megabytes (MB) or more of Synchronous Dynamic Random Access Memory (SDRAM), a smaller but faster memory cache that usually consists of about 512 Kilobytes (KB) of Static RAM (SRAM), and an even smaller cache “tag” that usually consists of about 16 KB to 64 KB of SRAM. The role of the memory cache is to provide, for some data requests, faster access to the requested data than the main memory can provide, and the role of the cache tag is to help determine whether or not the requested data is stored in the memory cache.\nIn one cache architecture, a microprocessor requests data from the memory system by first presenting the address of the requested data on a private cache bus interconnecting the microprocessor, the memory cache, and the cache tag. The cache tag receives the address (or, more commonly, a portion thereof), selects one of its internal memory locations in accordance with the address (or a portion thereof), and then writes out address data stored at the selected memory location to the microprocessor via the private cache bus. At the same time, the memory cache also receives the address, selects one of its internal memory locations in accordance with the address, and writes out the data stored at the selected memory location to the microprocessor.\nIf the microprocessor determines that the address data written out by the cache tag matches the address of the requested data (or a selected portion thereof), then a “cache hit” has occurred. In this circumstance, the microprocessor uses the data output by the memory cache, since the occurrence of a cache hit indicates that this data is the correct data. Conversely, if the microprocessor determines that the address data written out by the cache tag does not match the address of the requested data (or a selected portion thereof), then a “cache miss” has occurred. In this circumstance, the microprocessor requests the data from the main memory, because the occurrence of a cache miss indicates that the data output by the memory cache is not the correct data.\nDefects sometimes occur in a cache tag during the manufacturing process that prevent certain memory locations within the tag from functioning properly. If these defects cannot be repaired through conventional use of redundant elements, then the integrated circuit (IC) device that incorporates the defective tag is typically scrapped. There is, therefore, a need in the art for a device and method that can recover such devices for sale and subsequent use, thereby avoiding the need to scrap the devices."}
{"text": "Flying toys, especially kites, having long been used by all types of people for fun and entertainment. Kites come in all shapes and sizes. One particular type of kite is known as a rotatable airfoil kite. Such kites generally comprise a single elongated airfoil element or wing. This wing contains a circular disk member or stabilizer disk preferably attached at the center portion of the wing.\nVarious examples of rotatable kites are shown in U.S. Pat. No. 3,079,115 to Edwards, Jr., et al., in U.S. Pat. No. 4,012,017 to Springston, et al., in U.S. Pat. No. 4,121,794 to Lemelson, in U.S. Pat. No. 4,606,518 to Jeffrey, in U.S. Pat. No. 4,779,825 to Sams, and in U.S. Pat. No. 4,790,498 to Jeffrey.\nHowever, these known rotatable kites present several disadvantages to users. Generally, these kites are not very durable and tend to break on impact with the ground or other objects, or even due to a strong wind force. Also, these kites are expensive, cumbersome, hard to assemble, hard to transport, and unstable in flight. Further, these kites are incapable of being flown with tails or streamers.\nThe present invention solves the problems or disadvantages of known rotatable kites. The invention is lightweight yet very durable. It can be assembled with ease. The invention has a novel characteristic of being foldable for transportation or storage, while at the same time being stable in flight. Alternatively, an inflatable embodiment of the invention allows deflation and folding thereof for easy storage and transportation. The invention is specifically designed to be capable of flying with tails or streamers. Further, the invention is inexpensive and easy to fly.\nEdwards, Jr., et al. (U.S. Pat. No. 3,079,115) disclose a rotatable kite having a wing and a stabilizer disk. However, the kite is difficult to assemble due to plurality of tabs and slots which must be interconnected. Further, the kite is incapable of being folded or deflated for transportation and makes no provisions for flying with streamers.\nSpringston, et al. (U.S. Pat. No. 4,012,017) discloses a rotatable kite having a wing and two stabilizer disks. However, the kite is difficult to assemble due to the two disk design, and the need for bracing line. Further, the kite is incapable of being folded or deflated for transportation and makes no provisions for flying with streamers.\nLemelson (U.S. Pat. No. 4,121,794) discloses a relatively simple rotatable kite having a wing and a stabilizer disk. However, the kite is incapable of folding or deflation for transportation and makes no provisions for flying with streamers.\nJeffrey (U.S. Pat. No. 4,606,518) discloses another relatively simple rotatable kite having a wing and a stabilizer disk. However, the kite is incapable of folding or deflating for transportation and makes no provisions for flying with streamers.\nSams (U.S. Pat. No. 4,779,825) discloses various rotatable kites which are complex and difficult to assemble. Sams also discloses the use of cushioning strips covered with hinge tape to allow for movement of the stabilizer disk. However, these strips are used to prevent the disk from folding completely flat against the wing, unlike the present invention which is capable of such folding. Sams' disk is capable of movement during flight, unlike the present invention whose disk is stable during flight. Further, Sams makes no provisions for flying with streamers or deflation.\nJeffrey (U.S. Pat. No. 4,790,498) discloses a rotatable kite having a wing and a stabilizer disk. However, the kite is difficult to assemble due to a pair of closely spaced radially extending support elements. Further, the kite is incapable of folding or deflation and makes no provision for flying with streamers."}
{"text": "This invention relates to spreadable gelled food products and process for the production of the same.\nSyrups, such as honey, maple syrup, molasses, sugar cane syrups, and other viscous syrups principally containing sucrose, glucose, fructose and/or other sugars dissolved in water, find wide application as industrial food ingredients, in cooking, beverages, and as spreads. Syrups referred to here typically contain at least about 50%, and generally 70% or more, of dissolved sugars and are substantially self-preserving.\nOf all the syrups, honey has a particularly universal appeal. Whilst it has many uses as an industrial food ingredient and in cooking and beverages, honey is used most widely as a spread on bread, toast or crackers and the like.\nWhen used as a spread, honey has one much disliked aspectxe2x80x94that of being very runny or fluid. When spooned from a jar it is difficult to stop the honey from dribbling where it is not wanted. This aspect is even more unpopular when honey is applied to hot toast. The honey becomes even more fluid due to its absorption of heat and consequent reduction in viscosity. As a result, the honey dribbles on to the fingers and the clothes of the consumer, as well as on his plate or the table. Similar problems are found with other syrups.\nClearly, it is of advantage to be able to supply to consumers honey and other syrups in such a form that the above negative aspects of fluidity, especially at an elevated temperature, are overcome. Such an idea is not of itself new. Thus, it is traditional to convert honey to a candied form resulting in the formation of sugar crystals, preferably very fine crystals. The viscosity of candied honey is much higher than normal honey and therefore the tendency to run off toast or bread is reduced. However, candied honey (also referred to as creamed honey) often suffers from the opposite negative aspect in being so viscous that it is spread only with diffculty, often resulting in breakage of slices of bread or toast on which it is spread. This problem is a particular problem in cold weather, when candied honey becomes even more viscous.\nSeveral references exist which describe efforts to produce honey products which overcome the runny character of honey when used as a spread. Zelenak (Hungarian Patent T41974) describes a gelled honey requiring use of carrageenan, water, pH adjustment with citric acid, potassium sorbate as preservative, a buffer salt such as disodium hydrogen or sodium hexametaphosphate as well as some xanthan gum or sodium alginate. The final preferred soluble solids content of the gelled honey product was 65%. Thus the honey content in the final product was somewhat dilutedxe2x80x94no doubt as a consequence of the water required to dissolve the carrageenan and other gums and salts prior to mixing with honey.\nBrain et al (U.S. Pat. No. 4,532,143) produced a semi-gelled honey using a mixture of two different pectins preferably at a final soluble solids content of 62-68%, the final composition being substantially free of supplemental acid or alkaline earth metal cations.\nA minimum of 0.5% pectin was necessary to obtain the benefits of the invention the preferred range being 0.5-2%, more especially 0.75-1.25% of the pectins. The product finally was pasteurised. During processing the product is subject to high temperatures, such as from about 66xc2x0 C. to 94xc2x0 C. (150xc2x0 C.-200xc2x0 F.).\nIt is a characteristic of honeys that exposure to elevated temperatures results in deterioriation of flavour aroma and colour. The higher the temperature and the longer the exposure the more severe will be the quality deterioration. Temperatures in excess of about 60xc2x0 C. can be deleterious. A safer range is 50 to 55xc2x0 C. or even 45-46xc2x0 C. Quality aspects which can be lost as a consequence of excessive heating include loss of delicate flavour character, loss of brilliance and clarity and darkening of colour. Excessive heating can result in development of opacity to some degree. It can also result in the more rapid development of hydroxymethylfurfural which has a deleterious effect on honey quality.\nThis invention has as its main aim the production of a form of syrup, particularly honey, whereby the syrup loses to a substantial extent its tendency to be runny or fluid (especially at warm temperatures) without involving the mandatory use of added acids, salts, buffers or other mineral salts and not requiring the mandatory use of pectins, gums or other hydrocolloids or of mixtures of hydrocolloids other than agar-agar (agar). Surprisingly, it has been found that use of very low levels of agar, in amounts sufficient to form a gel, can be effective in changing the physical characteristics of the syrup as required. It is possible to produce a well-gelled form of honey or other syrup having a tender melting texture in the mouth but with a stable gel structure in storage even for extended periods as well as in a warm or hot situation as, for example, on hot toast. Thus, as little as 0.05% of agar will produce an acceptable result. From about 0.2 to 0.5% w/w, of agar in the final preparation is preferred. Agar may be used in excess of this indicated range if required. Because of the low concentration of hydrocolloid required, a relatively small amount of water or other aqueous solutions will suffice to pre-dissolve the agar before mixing with syrup. As a consesquence, a composition having a very high proportion of syrup, such as honey, results.\nIn accordance with a first aspect of this invention there is provided a spreadable gelled syrup composition which comprises said syrup, and agar in an amount sufficient to form a gel. The syrup is preferably honey. Some examples of other syrups which may be used in the invention include maple syrup, molasses, sugar cane syrup (sometimes known as xe2x80x9cgolden syrupxe2x80x9d and xe2x80x9ctreaclexe2x80x9d) and malt extract. The composition may additionally contain added water in which the agar is initially pre-dissolved.\nIn accordance with another aspect of the invention there is provided a process for the production of a spreadable gelled syrup composition which comprises dispersing solubilised agar through a syrup at a temperature which is above the gelation temperature of the agar, and thereafter cooling the composition so as to allow formation of a gel.\nIn accordance with a further aspect of this invention there is provided a process for the production of a spreadable gelled syrup composition which comprises dissolution of agar in water in the presence of heat, warming a syrup to a temperature which is above the gelation temperature of agar, and dispersing the dissolved agar uniformly in the warm syrup, and thereafter cooling the composition so as to allow formation of a gel.\nIn a further aspect there is provided a process for the production of a spreadable gelled honey composition which comprises dispersing solubilised agar through honey at a temperature which is above the gelation temperature of the agar and which is not deleterious to the honey, and thereafter cooling the composition so as to allow formation of a gel. On cooling to ambient temperature (between about 15xc2x0 C. to about 30xc2x0 C.) the composition remains fluid and non-gelled for about six hours or longer.\nIn a further aspect of the invention there is provided a process for the production of a spreadable gelled honey composition which comprises dissolution of agar in the presence of heat, warming honey to a temperature which is above the gelation temperature of agar and which is not deleterious to the honey, and adding the dissolved agar to the honey so as to disperse the agar therethrough, wherein on cooling to ambient temperature (between about 15xc2x0 C. to about 30xc2x0 C.) the resultant honey composition remains fluid and non-gelled for about 6 hours or more, after which gelation develops.\nIn contrast, according to the invention as presently claimed, a syrup, preferably honey, is heated to a temperature between 35xc2x0 C. and 60xc2x0 C. being a temperature non-deleterious to the syrup in particular honey solution, and above the gelation temperature of the agar (see claim 29).\nThe high concentration of honey in the final product may be achieved by the following features:\n1. Because a gelled product can be achieved by use of very low concentrations of hydrocolloid, very little water is required in which to pre-dissolve it. Hence the degree of dilution is very modest when this small amount of solution is mixed with honey. Thus the final honey concentration in the finished product can be in excess of 70% and typically in excess of 90%.\n2. By further processing as will be described hereafter, a gelled honey product can be prepared containing effectively no added water, so that the honey concentration in the finished product is about 99%xe2x80x94the only final non-honey component being the hydrocolloid used.\nAn advantage of using agar in the production of a spreadable, gelled honey composition is that the honey temperature throughout processing can be kept low, for example 50-60xc2x0 C. or even lower. The relatively small amount of hot agar solution has only slight impact on the temperature of the predominant proportion of honey in the composition. As agar has the unique property of not gelling below about 35xc2x0 C. to about 40xc2x0 C., the honey-agar blend can be held satisfactorily during processing and filling operations at moderate temperatures without the problems of pre-gelation as can occur with hydrocolloids such as pectins, alginates and the like. It has however further been found most surprisingly, that when agar is used in honey as will be described hereafter, the development of the gelled condition is delayed such that when the honey-agar mixture reaches ambient temperatures (about 15xc2x0 C. to about 30xc2x0 C., generally about 20xc2x0 C. to about 30xc2x0 C.) at which agar would normally gel, the mixture remains fluid. A tender spreadable gel begins to develop after about 6 hours or more. However, the firmness of the gel structure continues to develop over several days. This unexpected feature has important implications in relation to ease of industrial handling of a honey-agar gelled product, as will be described hereafter.\nIt is further observed that because a solution of agar in water is relatively fluidxe2x80x94much more fluid than honey and very much more so than other hydrocolloids or hydrocolloid mixturesxe2x80x94there is no increase and indeed actually a decrease in the tendency to air inclusion in the honey-agar product during blending and other industrial operations. Consequently deaeration of the final product is not required as a consequence of mixing the agar solution into the honey.\nThe gelled honey product according to the embodiment of the invention has a very tender easy to spoon and spread texture, with good clarity and a soft mouthfeel without any graininess. The honey flavour is readily released in the mouth so a full flavour impact is achieved. When the product is spread on hot toast it remains in its gel form and does not run off the surface of the hot toast as is the case with natural honey. The tendency of normal honey to dribble from a spoon is completely overcome. The gelled honey is tender, easy to spoon and spread and overcomes the spreading difficulties often associated with candied honey.\nThe invention is also applicable to spreadable gelled syrup products based on mixtures of 2 or more syrup-type foods with or without other flavourants or normal food components.\nSyrups which may be used in this invention include honey, maple syrup, molasses, sugar cane syrups, and other syrups which principally contain sucrose, glucose, fructose and/or other, sugars dissolved in water. Honey is particularly preferred and the following description refers largely to the use of honey. It is to be understood that the invention is not so limited and extends to compositions and processes involving other syrups including those discussed above, and combinations thereof.\nHoneys used in this invention may be of any type. There is, world-wide, an enormous variety of honeys. Some are very characteristic in their flavour, aroma and colour while others, while having generic honey characteristics, are not all that different or unique. It has been found that any type of honey is suitable for use in the products of this invention. Typically natural honeys vary from 80-82% soluble solids (SS) but some will be lower and others higher. Such variation in SS levels represents no problem in the implementation of this invention. Likewise the pH of honey, normally about 4, is quite suitable for production of a stable agar gel.\nAs mentioned above, agar has been found to give highly advantageous results in the production of spreadable gelled syrups, particularly honey. Agar is a hydrocolloid extracted from selected seaweeds and discovered originally in Japan. It is widely used in special foods, in microbiological media and also in special confectionery products. Various grades of agar are available and these differ in gel strength from as low as 150 up to 1000 g/cm2 or above (Japanese Agar Gel Strength Standard in xe2x80x9cFood Gelsxe2x80x9d Elsevier 1990, pp 23-25). It is preferred to use agar having a gel strengh from about 500 to about 1200 g/cm2. It is particularly preferred to use agar having a gel strength of about 1000 g/cm2, which is the grade of agar exemplified herein. It will be understood that use of grades of agar having different gel strengths will require adjustment of rates of usage.\nThe amount of agar used will depend upon the gel strength of the agar. Less agar is generally required in the situation where the gel strength is greater than 1000 g/cm2, compared with that required when the gel strength is less than about 800 g/cm2. For agars having a gel strength from about 500 to about 1200 g/cm2 a satisfactory gel is generally formed when the agar is present in an amount from about 0.05% to about 2.5% (w/w). Where the gel strength of the agar is, for example, between about 800 g/cm2 and 1200 g/cm2 it is generally preferred to utilise agar in an amount from about 0.05% to about 1.5% (w/w/).\nIn accordance with a first aspect of this invention there is provided a spreadable gelled syrup composition which comprises said syrup, and agar in an amount sufficient to form a gel. Where the composition contains added water, it is present in an amount from about 0.5% to about 15% (w/w). Very small amounts of agar may be used to form a gel with honey or other syrup. By way of example, agar in amount from 0.05% to 2.5% (w/w) may be used, particularly from 0.1% to 2%, more particularly from 0.2% to 0.5% w/w. The agar may be dissolved in a small amount of aqueous solution (generally water) in the presence of heat, as is described in more detail hereafter.\nThe spreadable gelled honey composition is in the form of a gel, this gelation being conferred by the agar. The honey forms part of the gel matrix. The gel is spreadable in the sense that it is capable of being spooned or spread onto an edible substrate. The gel is not fluid in the sense that it does not dribble like conventional honey, instead retaining an integral or coherent structure which is sufficiently soft to enable the gel material to be spread in a manner which is not dissimilar to pectin containing jams and like products.\nThe honey composition may comprise from 70% to about 99% weight/honey, the remainder comprising agar and water. Water may be used to dissolve the agar prior to its addition to the honey. This added water may be substantially removed from the honey composition utilising standard methods as described herein. As a consequence, the honey composition may comprise solely honey and agar. Agar is generally present in an amount from 0.05% to about 2.5%, such as from 0.1% to 1% w/w, more particularly from 0.1% to 0.5% w/w.\nThe soluble solids of the honey composition according to this invention may be in the range from about 70% to about 85%.\nIn one embodiment, the invention thus relates to a spreadable gelled honey composition which comprises from about 0.05% to 2.5% (w/w) agar, from about 70% to about 99% (w/w) honey, and optionally water. Added water (i.e. water other than that contained in honey or other syrup-type components), where present, may be in an amount from about 0.5% to about 15% (w/w). The composition preferably has a soluble solids content between about 70% to about 85% (w/w).\nHoney compositions which contain from 80% to about 99% honey, and soluble honey solids in the range from about 70% to 85% (w/w), which may be produced in accordance with this invention, have the distinct advantage of maintaining the colour, clarity and organoleptic properties of honey. The product is therefore essentially identical with regard to colour, taste and smell, with that of conventional honeys.\nSupplementary ingredients may be included in the honey compositions of this invention where it is desired to modify the taste or appearance of the honey. Supplementary ingredients may include spices or spice extracts, plant extracts, fruit flavours or fruit concentrates, flavour modifying agents such as bittering agents, salts, acids, and other commonly used food ingredients. When present such materials may be employed in an amount from about 0.01% to about 10% w/w of the final composition.\nHoney compositions may also include added sugars, such as sugars corresponding to those normally present in honey, namely glucose and fructose. Where other syrups are used, sugars which are conventionally found in such syrups may also be added. The addition of sugars, which may be dissolved in the agar solution prior to its addition to the honey or other syrup, have the effect of increasing the soluble solids content of the gelled product. Such materials may be included to increase the soluble solids content in an amount such as from about 0.5% to 10% w/w of the final composition.\nIn accordance with another aspect of the invention there is provided a process for the production of a spreadable gelled syrup composition which comprises dispersing solubilised agar through a syrup at a temperature which is above the gelation temperature of the agar, and thereafter cooling the composition so as to allow formation of a gel.\nIn another aspect this invention is directed to a process for the production of a spreadable gelled syrup composition which comprises dissolution of agar in water in the presence of heat, warming a syrup above the gelation temperature of agar, adding the dissolved agar to the warm syrup so as to disperse the dissolved agar therethrough, and thereafter cooling the composition so as to allow formation of a gel.\nIn a further aspect there is provided a process for the production of a spreadable gelled syrup composition which comprises dispersing solubilised agar through the syrup at a temperature which is above the gelation temperature of the agar and which is not deleterious to the syrup, and thereafter cooling the composition so as to allow formation of a gel. On cooling to ambient temperature (between about 15xc2x0 C. to about 30xc2x0 C.) the composition remains fluid and non-gelled for about six hours.\nIn a particularly preferred embodiment there is provided a process for the production of a spreadable gelled honey composition which comprises dissolution of agar in water in the presence of heat, warming honey to a temperature which is above the gelation temperature of agar and which is not deleterious to the honey, and adding the dissolved agar to the honey so as to disperse the agar therethrough, wherein on cooling to ambient temperature, which is below the gelation temperature of agar, the resultant honey composition remains fluid and non-gelled for a period of time thus facilitating processing of the honey composition. Generally, when cooled to an ambient temperature between about 15xc2x0 C. to about 30xc2x0 C. the honey composition does not gel for about six hours. At a temperature around 20xc2x0 C. gelation generally occurs at around 18 to 24 hours, a tender gel structure being formed. The firmness of this gel increases after about a further 24 hours.\nWhere it is indicated that the solution of agar in water is added to syrup which has been heated to a temperature which is above the gelation temperature of the agar solution, the underlying reason is to ensure that none of the aqueous agar solution is cooled to below a temperature at which pre-gelation of said aqueous solution may occur. In fact, once the aqueous agar solution is thoroughly admixed with the syrup its tendency to pre-gel is completely overcome, such that the syrup-aqueous agar solution, when thoroughly admixed, will not be subject to gelation for several hours after the mix cools to a temperature at which formation of a gel may occur.\nWhile it is not always practicable, it is conceivable that if the mixture of syrup and aqueous agar solution can be instantly mixed homogeneously, the temperature of the syrup would not have to be above the temperature at which the aqueous agar solution would gel.\nHowever, in practice, it is most convenient to warm the syrup to a temperature above the aqueous agar solution geling temperature. Such warming also reduces the viscosity of syrup and this improves its handling and mixing characteristics.\nA feature of an embodiment of this invention is that the honey is not subject to heating to the extent that its intrinsic qualities are negatively affected, that is, honey may be heated to a temperature which is above the gelation temperature of agar and which is not deleterious to the honey. It is common practice in the honey industry for honey to be carefully warmed to reduce its viscosity to allow for more efficient pumping, filtration, blending, deaeration and filling of the product into the final container. However, persons skilled in the art are careful not to exceed about 50-55xc2x0 C. and preferably 45-46xc2x0 C. to avoid deterioration of flavour, colour and brightness. At the same time, for some honeys it is necessary to heat the honey to about 60xc2x0 C. in order to avoid granulation (sugar crystallisation) on subsequent storage. In practising this invention it is not ncessary to exceed a product temperature of 50-55xc2x0 C. or even the most preferred temperature 45-46xc2x0 C. The fluidity of the product during preparation is such that it may be readily pumpted, filled and generally handled without problems.\nTo prepare a gelled honey product in accordance with an embodiment of this invention it is convenient, though not mandatory, for the honey to be carefully warmed to about 50xc2x0 C. in order to increase its fluidity and its ease of mixing. This is because a concentrated agar solution in water is required to be thoroughly and uniformly admixed throughout the honey. However, if a honey is especially heat sensitive it will suffice to warm the honey to about say 40-45xc2x0 C. though the viscosity at 40-45xc2x0 C. will be greater than 50xc2x0 C.\nThe gelation temperature of agar in water is somewhat variable depending on the grade, source and concentration of the agar. Generally, however, the gelation temperature of agar is normally 30xc2x0 C. to about 40xc2x0 C.\nSeparately a solution of agar in aqueous solution, generally water, is prepared. The amount of agar used is that which is capable of forming a gel in the presence of honey or another syrup or mixtures thereof. From about 0.05% to 2.5% (w/w) may be used, such as from 0.1% to 1% (w/w), more preferably 0.1% to 0.5% (w/w) of the final composition. A greater amount of agar may be used, however, it is preferred to use small amounts of agar such as from 0.01% to 1% (w/w). To achieve complete solution of agar in water, it is generally necessary to hold the agar-water mix at close to 100xc2x0 C. (or above) until all the agar particles are totally dissolved. Depending on the grade of agar used, solution can require as little as a few minutes or as much as half an hour or more. It is desirable to maintain a uniform proportion of agar and water so that the exact composition of the ultimate agar-water-honey mix is known. Hence during dissolution of agar evaporation of water is minimised. Alternatively, the proportion of agar to water is maintained by controlled addition of water to make up the water lost by evaporation.\nA solution of agar in water is relatively fluid, indeed much more fluid than is the case with most other hydrocolloids. This contributes to the ready dispersiblity of an agar solution throughout the honey, and also tends to reduce viscosity after mixing thoroughly with the honey.\nIn its simplest form a gelled honey product based on this invention comprises mixing a quantity of agar, completely dissolved in water, with a measured quantity of honey, generally in an amount of 70% to 95% w/w, at a temperature which will facilitate easy mixing without loss of honey quality. The honey and the agar solution may be mixed together at say 40-50xc2x0 C. (the agar may be cooled to 40xc2x0-50xc2x0 C., or added at a temperature from 50xc2x0 to 100xc2x0 C., given the small volume of addition), without any problems of agar pre-gelation or of overheating effects on honey. The mixture so produced will not gel immediately on cooling to a temperature at which a simple agar solution in water would gel, which, while somewhat variable depending on grade, source, concentration and other aspects, is normally about 30xc2x0 to 40xc2x0 C.\nSurprisingly, the honey-agar mixture does not set immediately on cooling to a temperature such as 15xc2x0 to 30xc2x0 C. which is below the normal gelation temperature of agar. This has a number of advantages. Firstly, there is no need to maintain a critical temperature for filling to avoid pre-gelation, as is the case with pectin gels and gels produced by other hydrocolloids or hydrocolloid mixtures. In industrial operations this is very important. For example, if a pectin-based product is allowed to pre-gel, the whole batch of product is lost and this is of considerable financial concern. Secondly, it is much easier to clean up industrial equipment if the remaining food has not formed a gel in the pumps, pipes, valves and other machine parts. Thirdly, since the basic product of the invention comprises honey (or other syrup) and agar solution, it is easy to standardise products awaiting final filling, the only test required being to confirm a desired level of soluble solids. By contrast, pectin gels usually require tests, in addition to soluble solids, such as pH, and setting capability, with other systems additionally require control of mandatory additions of acids, mineral salts and the like.\nThe above described mixing of an agar solution with honey results in a product which does not gel for at least about 6 hours at a temperature between about 15xc2x0 C. and 30xc2x0 C.\nBy about 24 hours the composition will show the initial formation of a gelled structure. By this is meant that the product, when spooned from a jar and placed on a flat surface. will to a large extent retain its shape and if inclined will not flow in a manner typical of normal honey. With the passage of time the firmness of the gel will increase such that after a few days of storage a gel having tender mouthfeel and excellent flavour release will have formed.\nAs mentioned above, additional sugars may be added to the honey composition in order to increase the SS of the final product and to off-set any dilution due to the addition of water in which the agar is dissolved. If it is desired to do this it is preferable to add the principal sugars which are naturally present in honeys. These sugars are fructose and glucose, and may be added in a proportion to each other which would not be conducive to granulation (crystallisation) of the honey sugars (for example 1.65:1 or more being the ratio of fructose to glucose, say 6.3 parts of fructose to 3.7 parts of glucose). This results in maintenance of a natural level of sweetness in the final product.\nAs a further alternative to the basic gelled honey product, a final product can be achieved containing only honey and agar, wherein the quantity of water required to dissolve agar can be first totally or almost totally removed from the base honey before mixing with the agar solution. Alternately the mixture of honey/agar solution can be treated to remove all or almost all of the mosture added as part of the agar solution. Removal of moisture as herein described can be effected at a comparatively low temperature (for example 50xc2x0 C.) by vacuum concentration until the required moisture removal is achieved. In the above procedures, that water required to pre-dissolve agar may be regarded as xe2x80x9ctemporaryxe2x80x9d being utilised solely for the purpose of achieving an agar solution, then being removed in effect after transferring the agar in a dissolved form into the liquid honey.\nThis invention is based on the surprising finding that agar may be used to produce a spreadable gelled honey product, which product maintains the colour, clarity and organoleptic properties of honey. Whilst the prior art has contemplated the use of carrageenans (Hungarian Patent T41974) and pectins (U.S. Pat. No. 4,532.143) to produce gelled honey, such proposals having been made in excess of ten years ago, it has hitherto never been contemplated to utilise agar in the production of spreadable gelled honey. The use of agar may not have been contemplated due to perceived problems of gel formation, adverse effects on colour, clarity and/or organoleptic properties of honey, or other factors."}
{"text": "1. Field of the Invention\nThe present invention relates generally to devices, methods, and systems of coded images. More particularly, an aspect of the present invention relates to lensless imaging and to spatially coded imaging.\n2. Brief Description of Prior Developments\nFixed (non-programmable) coded-aperture technology and programmable masks realized with Spatial Light Modulators were heretofore known in the art. The SLM elements are controlled so that the modulated transmission forms an aperture mask, encoding the image with a programmable function. The goal of decoding is to use knowledge of the aperture mask programmed with the programmable function to recover the object scene from the encoded image.\nInformation-theoretic analyses had identified two such classes of programmable masks used in association with recovering an object scene from an encoded image. The two classes are Uniformly Redundant Arrays (URAs) and Modified URAs (MURAs). Each may be used during the decoding process, however each has some drawbacks. A need therefore exists for an improved method for spatial decoding."}
{"text": "1. Field of the Invention\nThe present invention relates to an ESD protection device, and particularly to a low capacitance ESD protection device and a manufacturing method thereof applicable to RF circuits.\n2. Description of the Prior Art\nIn recent years, extensive research has been carried out on the possibility of using pure CMOS, rather than bipolar or BiCMOS technologies, for RF applications. It is forecasted that by 2003, a large number of applications will use CMOS technology for RF system-on-chip, such as remote control, radio modems, home system automation, etc. For example, in the home system automation, the concept is to replace the data cable between appliances by wireless link centered at 2.4 GHz, with a range of approximately 8 m. The reliability of the components used in these applications is therefore obviously very important.\nOne of the major bottlenecks for the transition from bipolar/BiCMOS to pure CMOS is the immunity of the circuits against Electro-Static Discharge (ESD). The detrimental effects of ESD in sub-micron devices have been widely known and reported. Further, the use of pure CMOS imposes a greater demand for adequate protection circuitry against ESD, due to the greater sensitivity of CMOS for electrical overstress in general and ESD in particular. RF circuits typically require a high performance and have tight design specifications. The conventional approach to achieve higher device ESD robustness is to incorporate on-chip ESD protection networks, as well as, in many cases external protection elements. As shown in FIG. 1, the standard approach to on-chip ESD protection involves incorporating protection at inputs, outputs, which consist of elements such as diodes, nMOSFETs, etc., provide a known current path to the ESD pulse away from the core circuitry. In the classical sub-micron CMOS technology devices, this common approach could yield enough ESD robustness if the protection circuits are implemented properly.\nFIGS. 1xcx9c3 are diagrams respectively showing three conventional ESD protection devices for CMOS, which are disclosed in xe2x80x9cESD Reliability Issues in RF CMOS Circuitsxe2x80x9d, M. K. Radhakrishnam, V. Vassilev, B. Keppens, V. De Heyn, M. Natarajan, and G. Groeseneken.\nFIG. 1 shows a cross-section of a ggnMOS (gate grounded nMOS) used for ESD protection. It includes a p substrate 11, isolation layer 12, n source and drain regions 131 and 132, p doped region 14, and a gate 15. The p doped region 14 couples the p substrate 11 to a cathode. The source region 131 and gate 14 are also coupled to the cathode. The drain region 132 is coupled to an anode.\nFIG. 2 show a cross-section of an oxide-isolated diode used for ESD protection. It includes a p substrate 21, oxide layer 22, a n or p well 23, n doped region 241 and p doped region 242 in the well 23, and poly-silicon layer 25. The p doped region 242 and n doped region 241 are respectively the cathode and anode of the diode.\nFIG. 3 shows a cross-section of a LVTSCR (low voltage triggered silicon controlled rectifier) used for ESD protection. It includes a p substrate 31, isolation layer 32, a transistor formed by an n source region 331, n drain region 332 and gate 34, p doped region 35, a n well 36, and n doped region 371 and p doped region 372 in the n well 36. The gate 34, source region 331 and p doped region 35 are coupled to the cathode while the n doped region 371 and n doped region 372 are coupled to the anode.\nThe previously described ESD protection devices provide ESD path between their cathode and anode. The cathode and anode may be coupled to power lines or I/O pads.\nHowever, for RF applications, the parasitic capacitance and the series resistance of the ESD protection structure directly impacts the RF performance. For example, an unaccounted parasitic capacitance of 1 pF from an ESD protection network or device at 2.4 GHz corresponds to a load of 66 Ohms. It is important to note that the typical input pin capacitance specification for a Low Noise Amplifier used in the 1-5 GHz is less than 100 fF.\nObviously, it degrades the performance of the RF circuits if the conventional ESD protection devices for pure CMOS are applied.\nThe object of the present invention is to provide a low capacitance ESD protection device and a manufacturing method thereof applicable to RF circuits.\nThe present invention provides a low capacitance ESD protection device. The device comprises a substrate, a well of a first conductivity type in the substrate, a first and second transistor of the first conductivity type respectively on two sides of the well, a guard ring of a second conductivity type in the substrate, surrounding the well, and the first and second transistor, and a doped region of the second conductivity type in the well, wherein profiles of a drain and source region of each of the first and second transistor are un-symmetrical.\nThe present invention further provides a method for manufacturing a low capacitance ESD protection device. The method comprises the steps of providing a substrate, forming a well of a first conductivity type in the substrate, forming a first and second transistor of the first conductivity type respectively on two sides of the well, wherein profiles of a drain and source region of each of the first and second transistor are un-symmetrical, forming a guard ring of a second conductivity type in the substrate, surrounding the well, and the first and second transistor, and forming a doped region of the second conductivity type in the well.\nThe present invention provides another low capacitance ESD protection device. The device comprises a substrate, a well of a first conductivity type in the substrate, a first and second transistor of the first conductivity type respectively on two sides of the well, a guard ring of a second conductivity type in the substrate, surrounding the well, and the first and second transistor, and a doped region of the second conductivity type in the well, wherein, for each of the first and second transistor, a width of a drain region is substantially equal to that of a source region and a length of the drain region is shorter than that of the source region.\nThe present invention also provides an integrated circuit comprising a core circuit, and a low capacitance ESD protection device protecting the core circuit from ESD damages. The low capacitance ESD protection device comprises a substrate, a well of a first conductivity type in the substrate, a first and second transistor of the first conductivity type respectively on two sides of the well, a guard ring of a second conductivity type in the substrate, surrounding the well, and the first and second transistor, and a doped region of the second conductivity type in the well, wherein profiles of a drain and source region of each of the first and second transistor are un-symmetrical."}
{"text": "1. Field\nThese inventions relate to lining fasteners, including for example insulation blanket fasteners, insulation blanket fasteners for aircraft, and including for example apparatus and methods of making and using lining fasteners.\n2. Related Art\nIn many aircraft, body portions such as the fuselage, passenger areas and/or cargo areas are lined with insulation panels and/or blankets. The panels are held in place with a plurality of fastener assemblies that are riveted, adhered or otherwise fixed to a fuselage structure such as a beam, stringers, frames or the like. Example fastener assemblies are molded from nylon or other plastic materials, and typically include a disc and a fastener stud having a base portion and a barb, where the barb may be used to puncture the blanket, for example. The barb typically includes a plurality of circumferential ridges or catches for engaging an opening in the disc placed over the barb. The plurality of ridges allows the disc to be placed at one of a number of axial positions along the barb as a function of the insulation panel thickness. For example, pressing the disc onto a barb past a given ridge on the barb may serve to hold the insulation panel in place, but further pressing the disc against the insulation panel may move the disc further along the barb and under another, further ridge for holding the insulation panel in place.\nThe fastener stud can be supported from the base portion on a portion of the fuselage, on a stringer, on a frame stud, or the like. The configuration of the base portion is generally formed to accommodate the profile of the underlying support structure."}
{"text": "The majority if not all houses require the construction of a bricked-in-gable under the roof. In constructing such gables, it is common practice to use a guide line suspended above the bricks that are being mortared into place under the gable. By following the guide line, the brick mason is assured that the stacked bricks and mortar remain even and level as the stack is built up even with the guide line suspended under the gable. Over the years, many devices have been utilized to suspend this guide line. Such a device is commonly referred to as a mason's guide line holder. In a classic arrangement for such a device, a pair of clamps suitable for holding a suspended line are attached to opposite ends of the gable face boards under the roof.\nOver the years, various improvements have been made to the classic arrangement with respect to having to un-fasten and move, and re-set and re-fasten the pair of clamps as the stack of bricks becomes taller.\nFor example, U.S. Pat. No. 3,571,931 (Williams) uses a pair of clamps suitable for holding a mason's guide line and can be affixed to the gable face boards on the opposite sides of the roof. The problem is that the clamps must be adjusted and moved over and over and they are attached at the gable face boards farthest apart. Constructing a bricked-in-gable under a large roof overhang becomes quite time consuming and labor intensive because of the all the un-fastening, re-fastening, and adjustment of the clamps and guide line.\nObviously, having to move the clamps and re-adjust and re-fasten them takes time and creates an inefficiency in the task of building a bricked-in-gable. U.S. Pat. No. 5,964,042 (Carper) shows a device for holding a mason's guide line that utilizes a pair of clamps and a pair of angle irons. The angle irons are attached to the underside of the gable and the clamps are slid up and down the angle irons and provide an alternative clamping surface and trackway. Unlike the present invention, when the gable overhang is greater than 4.5 inches, this device requires that angle irons be affixed to the underside of the gable by screwing or nailing the angle irons under the gable, thus adding additional labor to the task of building a bricked-in-gable. Also, additional expense is added to the manufacture and ultimate purchase price of the device because additional metal and fabrication must take place to produce the appropriate angle irons to be used with the clamps.\nAccordingly, the present invention describes an adjustable gable end line puller that solves all of the above problems."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a rack for carrying a snowboard. Specifically, the present invention relates to a snowboard rack, for holding and carrying a snowboard, designed to be mounted to a frame of a snowmobile.\n2. Background and Prior Art\nThe use of various types of equipment racks for carrying skis, wakeboards and snowboards are well known. More particularly, however, an example of a rack designed for carrying a snowboard in connection with operating a snowmobile is absent from the prior art. Commonly, snowboard racks in the prior art are taught to be attached to trucks, cars, or various types of water craft. As a result, the present invention seeks to provide a snowboard rack designed for mounting to a frame of a snowmobile for carrying a snowboard without hindering the snowmobile operator's ability to safely drive the snowmobile, and for securing the snowboard to the snowmobile in a safe manner to withstand various driving styles and conditions.\nThe following United States patents are herein incorporated by reference for their supporting teachings:                1) U.S. Pat. No. 5,752,638 issued to Meeks discloses a rack for storing one or more pairs of water skis and wake boards including one or more pairs of forks having parallel prongs extending from and coplanar with a base plate section. The pairs of forks are mounted on a support bar extending through an opening in the base plate section and may be positioned axially along and orientationally around the bar then clamped in the selected position such that each fork is aligned with its pair member. A strap is provided for each fork so that a water ski or wakeboard may be nested between neighboring prongs;        2) U.S. Pat. No. 6,019,265 issued to Deloza discloses a snowboard rack for transporting a snowboard in a pickup truck having a cab and a bed. The carrier has a carrier sleeve which extends over a header shaft to allow slidable movement thereupon. The carrier has a pair of cradles which support the snowboard therein, and a strap mechanism which secures the snowboard;        3) U.S. Pat. No. 6,189,753 issued to Kalhok et al. discloses a carrier rack for use primarily on a watercraft. The carrier rack has preferably two resilient, elastomeric cradles adapted to receive an object to be carried. Th cradles are adapted to be mounted on an outside surface of the watercraft, most preferably the gunwales;        4) U.S. Pat. No. 5,906,304 issued to Baldacchino discloses a rack that adjustably attaches to an inboard wall panel of a boat for securing water sport equipment in the cockpit of the boat. The device is configured so that the rack adjustably attaches to the inboard wall panel and secures a variety of equipment such as water skis and wake boards;        5) U.S. Pat. No. 5,390,837 issued to Ruffolo Jr. discloses a snowboard rack which is readily mountable to any roll-up window on an automotive vehicle; and        6) U.S. Pat. No. 5,564,683 issued to Stuck discloses a rotatable apparatus for servicing or repair of snowmobiles, including a frame or rack with a pair of generally parallel curved frame members which are located along a side of the snowmobile.        "}
{"text": "As shown in FIG. 1 (the drawing showing the various parts of the eye), the lens 11 is perhaps the most critical element within the eye in providing vision for the human or animal. It is suspended behind the cornea 12, the anterior chamber 13 and the iris 14 by the zonular fibers which connect it with the ciliary body of muscles 15 around its periphery.\nThe lens is composed of a central region, the nucleus 16, surrounded by a softer outer region, the cortex 17. It is encapsulated in the lens capsule 18o which is a thin transparent membrane. The front part of the membrane is called the anterior capsule 20, and the rear part, the posterior capsule 21.\nThe function of the lens is to focus light rays upon the sensitive retina 19. To focus light from a distant object, the ciliary muscles 15 relax, thus tightening the zonular fibers and reducing the thickness of the lens to its minimal dimension. To focus light from a near object, the ciliary muscles 15 release the tension on the zonular fibers to increase the thickness of the lens, the lens assuming a near spherical shape. This change increases the lens refractive power to again obtain focusing of the light rays on the retina 19.\nThe lens consists of about 35% protein and 65% water. When the proteinaceous material hardens or becomes suffused with minerals, the lens becomes cloudy or opaque. A cloudy lens or an opaque, non-functional lens is a cataract. Most cataracts are not visible to the casual observer until they become dense enough to cause blindness.\nCataract surgery involves the removal of the \"diseased\" lens from the eye. Two principal procedures for lens extractions are practiced currently. Intracapsular surgery involves removal of the lens together with its surrounding capsule. Where the posterior capsule may be attached slightly to the vitreous membrane, extracapsular extraction of the lens is much preferred. Rupturing the vitreous membrane, a definite hazard in intracapsular surgery, has serious consequences, including vitreous loss, vitreous hemorrhage, retina detachment etc.\nThis invention relates to extracapsular extraction, the procedure considered the safest. However, even in extracapsular surgery where a probe or similar instrument is used to remove the lens, either as a complete lens or by first dissolving the cortex and then shattering or breaking down the nucleus into smaller bits or pieces and extracting the bits, the danger of the surgical instrument penetrating through the posterior capsule exists, even with the most careful surgeon. Vitreous loss or hemorrhage, retina detachment, etc. may result.\nOne of the currently used procedures for extracapsular extraction of cataracts is disclosed in U.S. Pat. No. 3,996,935. For this procedure, the anterior capsule is first ruptured and removed, followed by removal of the cortex and nucleus of the lens, leaving the posterior capsule intact.\nThe primary object of the present invention is to reduce the chances to substantially zero of penetrating the posterior capsule in performing cataract surgery while retaining the anterior capsule substantially intact.\nA further object is to provide an instrument or device for achieving a successful procedure for removing unwanted objects from within the body without injuring neighboring body tissue.\nA still further object is to provide a device that will not only almost insure success in cataract surgery, but provide an option to the surgeon for simplifying the procedure, using a smaller incision to retain the anterior capsule substantially intact and to use less complicated secondary instrumentation than had previously been necessary."}
{"text": "1. Field of the Invention\nThis invention relates to an assembly structured to maintain a secure engagement between a louver and an associated carrier assembly of a vertical blind assembly and includes a support plate connected to and/or defining an upper end of the louver. A retaining clip is removably disposed in clamping engagement with anyone of a plurality of different types of connecting portions associated with the carrier assembly such that attachment between the connecting portion and the support plate is securely maintained.\n2. Description of the Related Art\nUse of vertical blind assemblies to cover windows, sliding doors, etc. has gained wide popularity over the last few years. Conventionally structured vertical blind assemblies normally include a header having an elongated configuration of sufficient length to extend along the upper periphery of the window, door, etc. intended to be covered. The header includes an interior track extending along substantially the entire length thereof and a plurality of carrier assemblies, at least equal in number to the number of louvers, blinds, vanes, etc. associated with the vertical blind assembly. Each of the carrier assemblies include structural features which facilitate easy passage thereof, as well as the louver or blind attached thereto, along the length of the track. A connecting portion is movably connected to each of the carrier assemblies and extends downwardly therefrom to an exterior of the header and in supporting engagement with a corresponding one of louvers.\nThe structural and operative features of each carrier assembly and associated connecting portion is such as to facilitate movement thereof, and the louver supported thereby, along the length of the track of the header assembly, as set forth above. In addition, each of the carrier assemblies are connected to a “tilt rod” also extending along the length of the header assembly and rotatably connected to mechanical linkage of each of the carrier assemblies so as to facilitate selective turning or tilting, in a concurrent, simultaneous fashion, of each of the louvers. Therefore, in typical fashion the plurality of louvers may extend along any portion of the length of the header in overlying relation to the corresponding door, window, etc. and/or be concurrently disposed into any type of slanted or tilted orientation so as to regulate the amount of light passing through the portal. Admittedly, vertical blind assemblies of the type generally described above include distinguishing structural features depending, at least in part, on a particular application for which a vertical blind assembly is intended. In addition, the vertical blind industry is replete with structurally modified components designed to improve the efficiency of the various types of vertical blind assemblies conventionally known and/or commercially available. Efforts to improve vertical blind assemblies include structural variations in the carrier assembly, connecting portion, track configuration, tilt rod, and selectively operable controls for regulating the position and orientation of the plurality of louvers or blinds associated with the vertical blind assembly.\nBy way of example, the structure of the various known carrier assemblies and their corresponding connecting portions may differ significantly. Common to a majority of such carrier assemblies and connecting portions is the ability of the connecting portion to be removably secured to an upper supporting end of a corresponding one of the louvers to which it is intended to be connected. Removable attachment between the connecting portion and the louver is provided to facilitate an effective connection there between while allowing quick and easy removal of individual ones of the louvers for repair or replacement. However, one disadvantage associated with the conventional manner of interconnecting the upper end of the louver to the connecting portion is the tendency for the louvers to become easily detached therefrom. Such undesirable detachment is particularly prevalent when the collection of louvers or blinds are forced along the length of the track and header assembly in a brisk manner, while a user simultaneously attempts to change the tilted or slanted orientation thereof.\nOther factors commonly affecting the stability of the interconnection between the connecting portion and the individual louvers is the weight, configuration and overall structure of the louver. By way of example only, a louver may be formed, at least in part, from heavy material components such as decorative chain links disposed in depending relation from and supported by a support structure disclosed adjacent an upper most end of the louver. Accordingly, when the aforementioned support portion is removably secured to anyone of a plurality of different types of connecting portions, the weight of the louver or blind may be such as to facilitate its detachment from the connection portion. Such inadvertent detachment is also common when the louver is subjected to unusual forces, such as engagement with people or objects, exposure to wind gusts, or rapid and simultaneous movement and tilting of the louvers through operation of the controls of the vertical blind assembly.\nTherefore, there is a significant and long recognized need in the vertical blind industry for an assembly which securely and consistently maintains a supporting interconnection of the connecting portion and individual louvers or blinds associated with the vertical blind assembly. Such an improved retaining assembly should be structured to not interfere with the normal operation of the vertical blind assembly, especially in terms of the collective movement of the louvers or blinds relative to the header assembly and track structure associated therewith. In addition, an improved retainer assembly of the type needed to overcome known disadvantages and problems of the type set forth above should have sufficient structural and operative versatility to be used with any of the different types of connecting structures known or commercially available, especially when such a large number of structurally distinguishable connecting structures are prevalently used in the vertical blind industry."}
{"text": "The present disclosure relates to a gas turbine engine, and more particularly to a reduction in purge air.\nThe core engine of a gas turbine engine typically includes a multistage axial compressor, a combustor and a high pressure turbine nozzle with one or more stages. Typical turbine nozzles, such as high pressure and low pressure turbine nozzles, define annular rings located adjacent to each turbine blade row to define axially alternate annular arrays of stator vanes and rotor blades.\nTo ensure that the rotatable blades and the static vane components do not contact each other under normal operating conditions, an annular gap is provided between the stator vanes and the bladed rotor. This requires, however, that the hot gases which pass through the turbine do not leak through the annular gap. Such leakage may result in a loss in turbine efficiency.\nThe conventional method to minimize hot gas leakage is the supply of high pressure purge air into the gap between the stator vanes and the bladed rotor. The purge air is directed radially outwardly over the surface of the rotatable disc and adjacent vane platform structure to exhaust through the gap into the core gas path. This minimizes hot gasses entrance into under-platform regions. These purge flows may cause some aerodynamic losses."}
{"text": "The present invention relates to a new and distinct cultivar of Buddleja plant botanically known as Buddleja hybrid and hereinafter referred to by the cultivar name ‘Bosedaze’.\nThe new cultivar originated in a controlled breeding program in Guadalupe, Calif. during September 2011. The objective of the breeding program was the development of Buddleja having distinctive flower coloration, good repeat blooming, and well-branched habits.\nThe new Buddleja cultivar is the result of cross-pollination. The female (seed) parent of the new cultivar is ‘Attraction’, not patented, characterized by its dark red-purple colored flowers, medium green-colored foliage, and moderately vigorous, upright growth habit. The male (pollen) parent of the new cultivar is ‘Miss Ruby’, U.S. Plant Pat. No. 19,950, characterized by its dark red-purple colored flowers, medium green-colored foliage, and moderately vigorous, compact-upright growth habit. The new cultivar was discovered and selected as a single flowering plant within the progeny of the above stated cross-pollination during August 2012 in a controlled environment in Guadalupe, Calif.\nAsexual reproduction of the new cultivar by terminal stem cuttings since August 2012 in Guadalupe, Calif. has demonstrated that the new cultivar reproduces true to type with all of the characteristics, as herein described, firmly fixed and retained through successive generations of such asexual propagation."}
{"text": "In a typical alarm system within a building, such as a fire or burglar alarm system, many types of sensors, detectors, lights, strobes, sounders and other associated devices may be located throughout the building as part of the system. Groups of these devices are often wired together along a pair of lines serving the communication and power requirements of the devices. A group of such devices on a pair of lines is often referred to as a \"line of devices\". Typically, many lines of devices can connect back to-a control panel that controls the overall operation of the alarm system.\nA typical line of devices, such as a line of smoke detectors, for example, may operate on 36 volts of line power supplied from a power source within or attached to the control panel.\nAccording to certain fire codes, alarm systems, such as those used in commercial establishments, must operate correctly during unintentional losses of AC power. To do so, most commercial alarm systems have battery backup power supplies. These alarm systems may also need to operate only on battery backup power during an emergency, for example, when AC power to the building is intentionally turned off. Typically, codes require alarm systems to be able to operate anywhere from twenty four to ninety hours before battery power is exhausted in a nonalarm condition.\nAn alarm system described in Takahashi, U.S. Pat. No. 5,151,683, provides intermittent alarm system device operation in order to conserve power when operating on batteries. In Takahashi, the control panel intermittently supplies power to each line of devices in a sequentially repeating order. While power is being applied to a line of devices, each device on that line is individually polled by the control panel for an alarm condition. In response to the poll, each device responds individually with a status message indicating either a normal or an alarm condition at that device. When all devices have been polled and have responded with a normal status, the application of power to the line of devices is terminated. The control panel then repeats the process with the next line of devices or with the same line after a short time interval has passed. In this manner, power is applied to the devices only when needed in order to check for an alarm condition at that device location. Thus, Takahashi conserves power by applying power to devices only when needed."}
{"text": "Whereas light sources are visible by their own emitted light, objects and materials appear to the eye according to how they modify incident light. The sensation of the color of an object is evoked by the physical stimulation of light-sensitive receptors in the human retina. The stimulation consists of electromagnetic radiation in the visible spectrum comprising wavelengths between about 380 and 780 nanometers.\nPerceived color of the object is the result of a combination of factors, such as: (1) the spectral power distribution of an illuminant emitted by a light source that is incident upon the object, (2) the modification of the spectral power distribution of the illuminant by the spectral reflectance or transmission characteristics of the illuminated object, (3) the excitation of light sensitive receptors in the eye by the modified light from the object, and (4) the perception and interpretation by the brain of signals produced by the light sensitive receptors.\nThe perception of color is attributed to the differing spectral sensitivities of the light sensitive receptors. The trichromacy of color sensation implies that many different spectral distributions can produce the same perceived color. Such equivalent stimuli, which produce the same perception even though they are physically different spectral distributions, are called metamers, and the phenomena metamerism. For example, it is known that the perceived color of an object can change quite markedly when the object is moved from incident daylight into incident artificial light. The spectrum of the illuminating light source is also known to have an effect on the perceived colors of a printed image in spite of the considerable physiological compensation that the eye makes for differences in illumination. Light sources of differing relative spectral power distributions are therefore known to have different color rendering properties; for example, light sources which emit very narrow band, or almost monochromatic, light are considered to render colors very poorly.\nAccording to the concept of metamerism, the respective colors of two objects may appear to be identical even though typically the spectral power distributions produced from the objects are different. Such power distributions, or stimuli, which are spectrally different but visually identical, are considered as metameric pairs. Because we measure light using only three cone types, the differences in these power distributions are indistinguishable. Two objects with different spectral reflectance functions may be perceived to match in color under one illuminant and not match under a different illuminant.\nCertain aspects of perceived color have been employed to disguise images by printing an image in one color and then overprinting the first image with a pattern in a different color having approximately the same apparent brightness. Adjacent zones of equal brightness appear to be visually blended, even though they are of differing colors, thereby confusing the perception of the original image.\nIt is known to print patterns in different colors such that the patterns may be viewed through one or more filters having certain correlated colors, such that the patterns will change, depending upon the colors involved. It is also known to print characters in different colors in an overlapping relationship such that the overlapped characters, when viewed through one colored filter, will give the appearance of only certain ones of the superimposed characters, and when viewed through a second and differing colored filter, will reveal certain other ones of the superimposed characters. Such approaches are known for encoding (or encrypting) information to prevent recognition of the information content of the pattern until the pattern is decoded and made comprehensible. These approaches have been applied to promotional gaming technology and in document security and document verification applications.\nTechniques are known for rendering flat, two-dimensional images that can stimulate an illusion of depth perception, that is, of a three-dimensional object or scene. Three-dimensional imaging can be classified into two major groups according to the quantity of information required to record the images: (1) binocular stereoscopic imaging, and (2) autostereoscopy, or three-dimensional spatial imaging. See Takanori Okoshi, Three-dimensional Imaging Techniques, Academic Press Inc., New York, (1976). Devices for performing binocular stereoscopic imaging include binocular viewers, parallax stereograms, lenticular-sheet binocular stereoscopic pictures, and binocular displays using Polaroid glasses or color filters. Devices for performing autostereoscopy include parallax panoramagrams, lenticular-sheet three-dimensional imaging, projection type three-dimensional displays, and integral photography.\nIn stereoscopy, a three-dimensional image is created by a series of two-dimensional images of an object captured from different perspectives, and therefore the three-dimensional image so produced contains multiple-angle information about the object. Physically displaced views of the same image are presented simultaneously to the eyes of an observer to convey the illusion of depth. These techniques typically employ a multiplexed pair of images, wherein the images are nearly identical and differ only so as to simulate parallax. The multiplexing is performed according to color, polarization, temporal, or position differences between the constituent images. For example, anaglyphic stereoscopy is a well-known process, in which left and right nearly-identical images are color-encoded by use of respective complementary color filters (e.g. cyan and red) for subsequent viewing through correspondingly colored lenses to separate the images as necessary for a simulated three-dimensional effect. When viewed through colored spectacles, the images merge to produce a stereoscopic sensation. The encoded image pair is known as an anaglyph, as it is typically rendered as two images of the same object taken from slightly different angles in two complementary colors.\nThis stereoscopic viewing of the multiplexed image pair typically requires the use of optical devices to channel each of the paired (left and right) images solely to the appropriate eye of the observer. A few autostereoscopic display techniques are known for providing subjectively three-dimensional viewing of a fixed image plane, without resort to eyewear and the like, by use of alternative devices based upon direction-multiplexed image displays. These devices typically employ optical diffraction, lenticular imaging, or holographic phenomena."}
{"text": "1. Field of the Invention\nThe present invention is related to business process management, and more particularly, to a unified architecture and a common syntactical framework for representing context of business process data, and management and tracking of the business processes.\n2. Description of the Related Art\nBusiness process management is an increasingly popular area of research and current software development. To date, there is no unified mechanism for a convenient representation of business processes data and data context. Typically, most software vendors provide one piece of the “puzzle”, such as CRM (customer relationship management), invoice tracking, order tracking, Gantt chart software, and so on. Many of the vendors in this area are relatively small companies, and typically the software product grows out of their own initial need to solve a particular problem, and the product is often tied to a particular world view of the vendor.\nFor a typical business enterprise, it is necessary to buy and integrate several pieces of software—for example, project management software typically does not “talk” to accounting software. Customer relationship management software does not talk to parts ordering software, and so on. Additionally, most such software packages have very limited ability to alter or change the way the data is presented to the user. Frequently, a relational database is used as part of the engine, with the attendant limited ability to manually add fields to the database. However, for most end users, the need to change the “candid” data representations represents a substantial difficulty, and often requires hiring consultants or additional IT staff.\nThe semantic web, see FIG. 1, comprises the standards and tools of XML, XML Schema, RDF, RDF Schema and OWL that are organized in the Semantic Web Stack. The OWL Web Ontology Language Overview describes the function and relationship of each of these components of the semantic web stack\nXML provides an elemental syntax for content structure within documents, yet associates no semantics with the meaning of the content contained within.\nXML Schema is a language for providing and restricting the structure and content of elements contained within XML documents.\nRDF is a simple language for expressing data models, which refer to objects (“resources”) and their relationships. An RDF-based model can be represented in XML syntax.\nRDF Schema extends RDF and is a vocabulary for describing properties and classes of RDF-based resources, with semantics for generalized-hierarchies of such properties and classes.\nOWL (see FIG. 2) illustrates additional vocabulary for describing properties and classes: among others, relations between classes (e.g. disjointness), cardinality (e.g. “exactly one”), equality, richer typing of properties, characteristics of properties (e.g., symmetry), and enumerated classes.\nAccordingly, what is needed is a unified architecture that permits a business enterprise to easily represent the context for business analysis in business processes that occur within the enterprise, and to provide an easy interaction mechanism between users representing various aspects of the business, and the data stored within the business model."}
{"text": "Advances in lightwave technology have made optical fiber a very popular medium for large bandwidth applications. In particular, optical technology is being utilized more and more in broadband systems wherein communications between systems take place on high-speed optical channels. As this trend continues to gain more and more momentum, the need for efficient utilization of the precious real estate on circuit boards, racks/shelves, back planes, distribution cabinets, etc., is becoming ever increasingly important. In order to fulfill expectations across the industry, opto-electronic modules and optic fiber devices need to continue to become miniaturized, thereby taking full advantage of the maturity of micro- and opto-electronic technologies for generating, transporting, managing and delivering broadband services to the ever increasing bandwidth demands of end users at increasingly lower costs. Thus, the industry has placed an emphasis on small form factor optical connectors, such as the LC connector from Lucent Technologies, Inc. However, miniaturization is tempered by the requirements of transmission efficiency. For instance, with the advent of new standards such as gigabit Ethernet, wherein the transmission efficiency is becoming more and more critical, the performance of optical connectors is becoming correspondingly important for healthy operation of the system. Thus, it is desirable to obtain component miniaturization without sacrificing transmission efficiency, and preferably while improving transmission efficiency.\nWith the miniaturization of optical modules and optical fiber devices, the management of optical fiber congestion has become an issue at optical interfaces and connection distribution points. One solution is the use of multi-fiber ribbon in which a plurality of optical fibers are organized and molded side by side in a plastic ribbon. It is known to interconnect these ribbon cables by supporting the fibers between two support members made of a monocrystalline material, such as silicon. In the support members are V-grooves formed utilizing photolithographic masking and etching techniques. The fibers are placed side by side in individual V-grooves of one support member and the other mating support member having corresponding V-grooves is placed over the fibers so as to bind or hold the fibers in a high precision, spatial relationship between the mating V-grooves. The top and bottom support members sandwiching the multi-fiber ribbon are typically bonded together with a clamp or adhesive, forming a ferrule of a multi-fiber connector. Two mating ferrules with the same fiber spacing may then be placed in an abutting relationship so that the ends of the fibers of the respective ferrules are substantially co-axially aligned with one another, thereby forming a multi-fiber connection. If desired, such ferrules can be stacked in order to increase the interconnection density.\nMulti-fiber ribbons and connectors have numerous applications in optic communication systems. For instance, some opto-electronic and optical application specific integrated circuits (OASIC) devices, e.g, optical switches, optical power splitters/combiners, routers, etc., have several input and/or output ports arranged as linear arrays to which a plurality of fibers are to be coupled. Further, since optical fibers are attached to launch optical signals into these devices and extract optical signals out of these devices, splicing of arrays of fibers (i.e., a multi-fiber ribbon) to such devices can be achieved using mutifiber connectors. Yet another possible application relates to an optical fan-out fabric where an array of fibers in a multi-fiber ribbon may be broken into simplex or duplex channels for distribution purposes, as is often desired.\nA critical factor to the optical efficiency of a multi-fiber connector, whether or not stacked, is the precise alignment of the mating ferrules with regard to one another. As the ferrule structures utilized to achieve the precise axial alignment of corresponding optical fibers of a multi-fiber connection become smaller, there exist a need for connectors that are likewise space efficient so that full advantage of the miniaturized ferrule can be realized with higher interconnection density. Further, there also exist the need for multi-fiber connectors to be user friendly so that the operation and utility of the multi-fiber connector is intuitive to the workers that will be installing systems utilizing optical components. For example, it is desirable for multi-fiber connectors to have plug-and-play capability, in that they can be quickly and easily coupled to a piece of equipment, device, or one another.\nIn summary, there continues to exist a need for miniaturized multi-fiber connectors in order to take advantage of the more space efficient optical ferrules while, at the same time, increasing the functionality and ease of use of such multi-fiber connectors."}
{"text": "In the service industry, delivery of services with expected quality is of utmost importance in order to fulfill end user requirements. In the present scenario, the global economy highlights issue of significantly increasing pressure experienced by the service industry to operate at higher speeds to deliver the services with the quality desired and value to cope up with cut throat competition present in the global market. An information technology (IT) service is one of among various other services which have always played a key role in dealing difficulties related to delivery of various services for the industries or organizations. The IT service delivered by an IT vendor may result in generation of a business value for the industry or the organization. Generally, the business value of any industry or an organization is articulated quantitatively based on Return on Investment (ROI) which is widely used to measure the effectiveness of the investment made for the IT services in the industry and also to compute the business value generated by the industry or the organization after implementing the IT services.\nMost of traditional methods of articulation of the business value are based on improvement in the quality levels of the service delivered by the IT services. One of such methods, a business value gained by an organization for the qualities procured in the IT services is determined in terms of value realization index (VRI). In this method, the VRI is articulated based upon delivered quality level and promised quality level of the IT services. In the method for articulation of the VRI, the articulation of the VRI is dependent on a single-scale of 0 to 100. Further the present method fails to represent various quality levels of the IT services and hence failed to represent direct measures of various qualities of the IT services."}
{"text": "1. Field of the Invention\nThe invention relates to audio peak limiters and other audio control systems, more particularly to the cancellation of lower sideband modulation distortion in such systems.\n2. Prior Art\nFast audio peak limiters include simple peak clippers, radio frequency clippers, Hilbert transform limiters as disclosed in U.S. Pat. No. 4,495,643, delay-line limiters, and others. Some of these operate directly on an input audio signal such as the simple peak clipper, which replaces any positive or negative instantaneous level exceeding a preset positive or negative threshold with the corresponding threshold value. Other peak limiters operate by multiplying a control signal by the input signal. All systems that operate directly on the signal have equivalent systems (producing the same outputs) that multiply the input signal by a control signal. U.S. Pat. No. 5,168,526 discusses how such an equivalent system can be synthesized.\nAny system in which the input signal is multiplied by a control signal can be considered an amplitude modulator, where there are multiple \"carriers\" corresponding to each Fourier frequency component of the input signal and the \"modulating signal\" is the control signal. It is well known that amplitude modulators produce a pair of symmetrical, mirror-image sidebands on each side of the carrier. The frequencies in the upper sideband are the sum of the carrier and each Fourier frequency component in the modulating signal. The frequencies in the lower sideband are the difference between the carrier and each Fourier frequency component in the modulating signal. In the case of the peak limiting systems, the modulating process adds to each Fourier frequency component of the input signal a pair of sidebands corresponding to the frequency components in the control signal. This is true of all prior art processes enumerated above.\nPsychoacoustic research has established that strong frequency components in an audio signal have the ability to \"mask\" weaker signals so they cannot be heard in the presence of the strong signals, although they would otherwise be audible in the absence of a strong signal. As the weak frequency becomes more distant in frequency from the strong signal, the strong signal must become stronger and stronger to mask the weak signal.\nThis effect is asymmetrical with frequency. Strong signals are more able to mask weak signals \"x\" octaves above a strong signal than to mask weak signals \"x\" octaves below a strong signal. Therefore, in any of the prior art peak limiting systems discussed, the lower sideband (representing distortion induced by the peak limiting process) is much more likely to be audible than is the upper sideband because the undistorted input signal is less able to psychoacoustically mask the lower sideband. This is particularly true because the lower sideband is arithmetically symmetrical to the upper sideband, occupying many more octaves than does the upper sideband."}
{"text": "This application is based on Japanese Patent Applications No. 8-349939 filed on Dec. 27, 1996 and No. 9-059600 filed on Mar. 13, 1997, the entire contents of which are incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to data communications technologies, and more particularly to real time data communications technologies. A xe2x80x9creal timexe2x80x9d response to an event is essentially simultaneous with the event itself. However, in communications, because of time delay for transmission time, signal synchronization, other necessary signal process or the like, xe2x80x9creal timexe2x80x9d does not mean strictly simultaneous.\n2. Description of the Related Art\nAs a standard specification for communications between electronic musical instruments, a music instrumental digital interface (MIDI) specification is known. Electronic musical instruments equipped with interfaces of the MIDI specification can communicate with each other by transferring MIDI data via a MIDI cable.\nFor example, an electronic musical instrument transmits MIDI data of a musical performance by a player, and another musical instrument receives it to reproduce it. As one electronic musical instrument is played, another electronic musical instrument can be played in real time.\nIn a communications network interconnecting a plurality of general computers, various types of data are transferred. For example, live musical tone data or other MIDI data can be transmitted from one computer, which once stored the data in its storage device such as a hard disk, via the communications network to another computer which stores the received data in its storage device. A general communications network is, however, configured to perform only general data communications, and is not configured to properly process MIDI data.\nSpecifically, although the MIDI specification allows the xe2x80x9creal timexe2x80x9d communications to be performed between electronic musical instruments, it is not suitable for long distance communications and communications via a number of nodes. The general communications network is essentially configured to provide services of long distance communications and multiple-node communications, but it does not take account of xe2x80x9creal timexe2x80x9d communications between electronic musical instruments.\nReal time communications of musical information uses a large amount of information per unit time, and the traffic of the communications line becomes heavy. As compared to point-to-point communications, point-to-multipoint communications of musical tone data is more likely to make the traffic of communications lines heavy. The heavy traffic of communications lines generates a transmission delay and hinders a real time musical performance.\nIt is an object of the present invention to provide technologies of musical tone data communications capable of a real time musical performance at multiple nodes.\nIt is another object of the present invention to provide technologies of data communications capable of avoiding a heavy traffic of communications lines.\nAccording to one aspect of the present invention, there is provided a musical tone data communications system, comprising: transmitting means for transmitting inputted MIDI data in real time over a communication network.\nAccording to another aspect of the present invention, there is provided a data communications system comprising: receiving means for receiving data; access checking means for checking the number of communications lines accessed externally; and transmitting means capable of reducing the amount of data received by the receiving means in accordance with the number of communications lines accessed externally, and transmitting the reduced data to the communications lines accessed externally.\nIf the number of accessed communications lines is large, the amount of received data is reduced to thereby alleviate the traffic congestion, whereas if the number of accessed communications lines is small, it is not always necessary to reduce the data amount.\nAccording to a further aspect of the present invention, there is provided a communication system having a plurality of communications apparatuses each having receiving means and transmitting means, wherein: the receiving means of the plurality of communications apparatuses receive the same data; the transmitting means of the plurality of communications apparatuses can reduce the amount of data received by the receiving means and can transmit the reduced data; and the data reduced by one of the communications apparatuses is different from the data reduced by another of the communications apparatuses.\nSince the data reduced by one and another of communications apparatuses is different, the quality of data transmitted from each communication apparatus is different. For example, the type or reduction factor of the reduced data may be made different at each communication apparatus. Therefore, a user can obtain data of a desired quality by accessing a proper communication apparatus.\nAccording to still another aspect of the invention, there is provided a musical tone data communications method comprising the steps of: (a) transmitting MIDI data over a communications network; and (b) receiving the transmitted MIDI data and supplying the received MIDI data to a tone generator in real time.\nMIDI data can be transmitted to a number of nodes by using a communications network. At each node, the MIDI data is reproduced in real time to generate musical tones.\nAccording to still another aspect of the invention, there is provided a musical tone data communications method comprising the steps of: (a) transmitting MIDI data; and (b) transmitting recovery data after the MIDI data is transmitted, the recovery data indicating a continuation of transmission of the MIDI data.\nIf there is no communications error, transmitted MIDI data can be correctly received at a partner communications apparatus. If there is a communications error, transmitted MIDI data cannot be correctly received at a partner communications apparatus. Even in such a case, the communication error can be remedied by transmitting the recovery data."}
{"text": "An assembly as well as systems and methods for performing parallel crystallisation experiments are known from WO 02/06802. The known assembly comprises a microplate having multiple microwells having an opening at the top. Sealing of the wells is effected by O-ring seals around the top of each well, which are interpositioned between the microplate and a cover plate.\nThe known assembly is not sufficiently practical when conducting parallel experiments in high volumes, known as high throughput experimentation. In particular the known assembly is unsatisfactory when conducting parallel crystallisation experiments."}
{"text": "The patterning of radiation sensitive polymeric films with high energy radiation such as photons, electrons, or ion beams is the principle means of defining high resolution circuitry found in semiconductor devices. The radiation sensitive films, often referred to as photoresists regardless of the radiation source, generally consist of multicomponent formulations that are coated onto a desired substrate such as a silicon wafer. The radiation is most commonly ultraviolet light at wavelengths of 436, 365, 257, 248, 193 or 157 nanometers (nm), or a beam of electrons or ions, or ‘soft’ x-ray radiation, also referred to as extreme ultraviolet (EUV) or x-rays. The radiation is exposed patternwise and induces a chemical transformation that renders the solubility of the exposed regions of the films different from that of the unexposed areas when the films are treated with an appropriate developer, usually a dilute, basic aqueous solution, such as aqueous tetramethylammonium hydroxide (TMAH).\nTypical photoresists contain a polymeric component and are generally comprised of a polymeric matrix, a radiation sensitive component, a casting solvent, and other performance enhancing additives. The highest performing photoresists in terms of sensitivity to radiation and resolution capability are “chemically-amplified” photoresists, allowing high resolution, high contrast and high sensitivity that are not generally provided by other photoresists. Chemically amplified photoresists are based on a catalytic mechanism that allows a relatively large number of chemical events such as, for example, deprotection reactions in the case of positive photoresists or crosslinking reactions in the case of negative tone photoresists, to be brought about by the application of a relatively low dose of radiation that induces formation of the catalyst, often a strong acid.\nAlthough chemically-amplified resists have been developed for 248, 193 and 157 nm lithography, certain barriers to achieving higher resolution and smaller feature sizes remain due to physical, processing and material limitations. One such phenomenon that arises for imaging in the sub-50 nm regime, resulting in diminished image integrity in the pattern, is referred to as “image blur” (see, e.g., Hinsberg et al., Proc. SPIE, (2000), 3999, 148 and Houle et al., J. Vac. Sci. Technol B, (2000), 18, 1874). Image blur is generally thought to result from two contributing factors: gradient-driven acid diffusion and reaction propagation, the result being a distortion in the developable image compared to the projected aerial image transferred onto the film. Although, both factors contribute to image blur, the degree of the effect from each is different. Temperature also has a differing effect on each factor.\nAcid diffusion is further thought to depend on several factors, including the type of photoacid generator (PAG) and the mobility in the photoresist polymer. In turn, the acid mobility in the polymer is dependent on a variety of factors, including, among others, the chemical functionality of the polymer, and the temperature and time of baking during resist processing.\nReaction propagation likewise depends on a number of factors, such as the activation energy (enthalpy) and the volatility of products (entropy).\nThus, as the need for better resolutions, minimum feature sizes, improved sensitivity and process latitude increases, image blur due to both contributing factors must be minimized. While both may be reduced to a degree by the use of acid-quenchers, or bases, the extent of thermally induced image blur estimated to be on the order of 10-50 nm with conventional resists and processing (see also Breyta et al., U.S. Pat. No. 6,227,546) suggests that improvements are necessary in order to achieve sub-50 nm imaging.\nAn ongoing need therefore exists for new photoresist materials and compositions, as well as methods of patterning substrates, which can lead to improved high resolution photoresist applications."}
{"text": "Billboards are structures typically used to display large-scale signage in outdoor locations. These billboards may be free-standing, such as those seen along highways and roads, or installed on another structure, e.g., atop a building or stadium. During daylight hours, the signage may be visible without any additional lighting. Traditional billboards usually contain lighting, such as flood lights, installed on the front of the billboard structure and directed toward the signage in order to light the sign and improve visibility, for instance, during nighttime hours. Newer types of billboards, such as LED (Light Emitting Diode) billboards, use a different lighting system and require an entirely different structure than the lighting system and structure of traditional billboards."}
{"text": "Inductors are known in the art, and are used to resist against changes in current through the coil. Inductors typically include a coil of conductive material wrapped around a magnetic core. Often, such cores are formed in a closed loop. Known inductors include coils that are wrapped manually, such as by a winding machine. Typically, the inductor's magnetic core and windings are placed between an outer wall and an inner wall.\nIn some applications, inductors dissipate significant quantities of heat. Because of this, known inductors are potted in heat dissipating materials. The pottant is typically poured between the inner wall and the outer wall to surround the windings and to provide environmental, thermal, and structural support to the cores and windings. Pottants must have a high degree of plasticity to fully fill the cavity between the windings and the outer casing when poured. Furthermore, the pottant selected should have as high of a coefficient of thermal transfer as possible, in order to maximize heat transfer to the outer casing.\nKnown pottants attempt to provide both desired rheological attributes (i.e., high plasticity/flowability for pouring) as well as high coefficients of thermal transfer."}
{"text": "The present invention relates to a differential amplifier having high gain at a low frequency signal input."}
{"text": "With the increasing clock speed of integrated circuits (ICs), the relative timing between two or more signals received at the same component is becoming more and more of a design consideration and constraint. Generally, in ICs, there is a typical propagation rate through a given component of an IC. Therefore, the time of propagation of a signal that arrives over a first signal path is synchronized with a computer clock, and a second signal that arrives over a signal path is also synchronized with the computer clock. These signal path propagation times can be calculated and appropriate design choices made. However, if two signals are received at the separate inputs at substantially the same time or relatively close in time to one another, the propagation speed of the resulting signal through that component can be either substantially shorter or substantially longer than the design considerations typically allow for. This can create two potential error conditions in chip design.\nThe first error condition is generally known as a “late mode” failure, wherein the signal takes longer to propagate, due to simultaneous reception by a component of two incoming signals, than was accounted for in the ideal “non-simultaneous reception” case. The “late mode” case can typically be compensated for slowing the computer clock cycle speed, thereby allowing the signal time to get to the next clock-synchronized component.\nThe second error condition is an “early mode” failure. In a race condition, due to the substantially simultaneous receipt of two different signals, the resulting signal propagates out of the component faster than anticipated. This is especially irksome, as typically the simultaneous receipt of two independent signals is itself not a function of a received clock cycle, so slowing down the clock cycle typically will not prevent an early condition. Conventional technologies and methods for calculating acceptable design margins to account for both the early and late conditions tend to be either unwieldy, unacceptably inaccurate, or both.\nTherefore, there is a need for a method and/or system for estimating design parameters for both early and late race conditions in IC design that addresses at least some of the problems associated with conventional methods and/or systems for estimating design parameters for both late and early race conditions."}
{"text": "“Health care professionals” is used broadly here to refer to anyone who participates in the diagnosis or treatment of medical problems. For example, medical doctors, dentists, nurses, nurse-practitioners, medical technologists, physical therapists, and other health workers that assist in examination of patients, diagnosis, or treatment are all included by this term.\nA health care professional diagnoses an illness by collecting and evaluating information about the patient, then determining what disease or condition best fits the information. The information gathered from the patient usually is processed to reach a diagnosis by using a protocol learned during the professional's professional training and modified and updated by his or her medical experience. The protocol is an ordered process by which a health care professional ascertains information that allows the professional to rule out possible diseases until enough information is gathered to eliminate all but the diagnosed condition. Alternatively, the protocol may end when an appropriate treatment is identified. Recently, medical associations, health maintenance organizations, and hospitals, among others, have prescribed protocols. Employed health care professionals in particular are often subject to mandated protocols.\nOne problem in the field of medicine is how to improve diagnostic protocols to take into account advances in medical knowledge. A related problem is how to ensure that health care professionals update their skills to take advantage of advances in medical knowledge. Still another problem is how to expedite the diagnosis and treatment of certain conditions that should be treated quickly, so treatment can begin soon enough to be most effective.\nU.S. Pat. No. 6,095,973 discloses a data processing system and method for evaluating the treatment of chest pain patients in a medical facility.\nU.S. Pat. No. 6,029,138 discloses a decision support system for the selection of a diagnostic test or therapeutic intervention, which are both called “studies” in that patent.\nU.S. Pat. No. 4,857,713 discloses a program for reducing hospital errors in the delivery of medications, goods, services or procedures in patient treatment.\nU.S. Pat. No. 5,732,397 describes an automated system for use in decision-making processes, which is said to improve the quality and consistency of decisions made.\nU.S. Pat. No. 5,772,585 discloses a common user interface to allow different medical personnel access to centralized files regarding patients.\nU.S. Pat. No. 5,832,450 describes an electronic medical record system that stores data about individual patient encounters in a convenient form.\nU.S. Pat. No. 5,845,255 describes an electronic prescription creation system for physician use that includes an adverse indication review and online access to comprehensive drug information including scientific literature.\nU.S. Pat. No. 5,911,132 discloses diagnosing and treating patient diseases using a epidemiological database containing medical, personal or epidemiological data relevant to a presented set of symptoms, test results, a diagnosis, etc.\nU.S. Pat. No. 5,915,240 discloses a context-sensitive medical lookup reference computer system for accessing medical information over a network.\nU.S. Pat. No. 5,924,074 discloses a medical records system that is said to create and maintain all patient data electronically.\nU.S. Pat. No. 5,953,704 discloses a system in which a user inputs information related to the health condition of an individual.\nU.S. Pat. No. 6,022,315 discloses a system and method for providing computerized, knowledge-based medical diagnostic and treatment advice to the general public over a telephone network or a computer network.\nThere is currently a need in the medical field for a system that communicates to a health care professional carrying out a diagnosis that a certain symptom, combination of symptoms, or other patient information recorded by the physician is associated with an increased risk of a missed medical care opportunity leading to a less favorable patient outcome. (A “medical care opportunity” is defined as an opportunity to correctly or more quickly diagnose or treat the patient's condition and thus provide a better patient outcome.) Further, there is also a need in the medical field for a system for communicating to the health care professional special steps to take to avoid the missed medical care opportunity."}
{"text": "Vehicle based computing systems, such as the FORD SYNC system are growing in popularity. Using various sources of vehicle information, driver inputs and connections to vehicle systems, the SYNC system can add a variety of functionality and novelty to the driving experience.\nFurthermore, systems such as SYNC can often communicate with remote devices either to gain information from those devices, or to use those devices to access a remote network. For example, in one instance, SYNC can communicate with a cellular phone, and use the cellular phone's ability to communicate with a remote network to send and receive information to and from the remote network. In another example, SYNC can query a GPS navigational device, such as a TOMTOM, and receive navigational information.\nIn addition to querying a device, such as a TOMTOM to receive navigational information, SYNC can also communicate with the TOMTOM and provide instructions, often comparable to pressing a selection on the TOMTOM's screen, through the SYNC system. The instructions can be provided, for example, by a spoken driver command processed through the SYNC system."}
{"text": "Prior to this invention various fluorometers have been used in a wide variety of applications for measuring the fluorescence of fluorescently excitable materials. For example, fluorometers are used in junction with fluorescent assays to detect and measure the quantities of immunological and non-immunological substances.\nIn carrying out fluorescent assays, microtest plates (or microtitration plates, as they are also called) and strips of microtest wells are often used. Microtest plates are formed with a multiplicity of wells which are joined together in a molded one-piece structure for containing microliter quantities of fluid samples in liquid or solid form. Examples of a microtest plate and microtest wells in strip form are described in U.S. Pat. No. 4,154,795 which issued to A. C. Thorne on May 15, 1979.\nThe use of microtest plates and microtest strips of wells in fluorescent and other types of assays offers several important advantages. First, they permit the mass preparation of a large number of test sample solutions at the same time. Second, they are more convenient to handle as compared with individual test tubes. Third, they can easily and inexpensively be washed. Fourth, they are inexpensive and disposable. Fifth, they are customarily formed from plastic materials which are not fragile like glass. Sixth, they can be made from a material having an ability to attract certain molecules such as protein molecules so that they can serve as a solid phase in an immunoassay.\nSome fluorometers are not sufficiently sensitive to measure the fluorescence of the small, microliter quantities of the relatively low fluorescent samples which are prepared with the microtest plate and strip equipment described above. Other fluorometers, while having sufficient sensititvity, are usually unsuitable for measuring the fluorescence of substances in microtest wells because they are designed to direct the exciting light and/or the sample's emitted light through a wall of the sample-holding vessel. As a result, the microtest plates and strips, which are customarily molded from plastic materials having a substantial level of native fluorescence, are excessively excited to produce spurious light emissions which interfere with and impair accurate measurements of the intensity of the light emitted by the fluorescently excited test sample itself."}
{"text": "The present invention relates to an apparatus for recycling the exhaust gas of the engine crankcase of a car which filtrates the exhaust gas of the engine crankcase permitting the filtrated gas to be further guided into the engine of the car for burning.\nDuring the operation of the engine of a car, fuel mixture may partially leak out of the combustion chamber to the engine crankcase through the gap between the piston ring and the cylinder wall of the engine. Leaked fuel mixture mixed with exhaust gas of the engine crankcase must be properly treated otherwise it will pollute the air. Various economizers have been disclosed to collect the exhaust gas of the engine crankcase for recycling. The positive crankcase ventilating system (PCV system) is one of the most commonly used systems for recycling the exhaust gas of the engine crankcase. However, the PCV system cannot completely recycle the exhaust gas of the engine crankcase and partial exhaust gas will come out of the engine crankcase to pollute the air. Beside this Carbon will be precipitated on the inside wall of the engine crankcase when engine gas is retained in the hot engine crankcase, causing the air and oil holes of the engine blocked or partially blocked, and therefore the service life of the engine will be greatly shortened. Furthermore, when the speed of the car is reducing or the car is stopped while the engine is operating at a high speed, the throttle valve is closed, causing a high negative pressure produced at the intake manifold of the engine. Under this stage, the suction force of the high negative pressure will stop the valve of the PCV system, causing the cylinder of the engine unable to obtain sufficient air for combustion, and therefore carbon monoxide will instantly increase."}
{"text": "The present disclosure relates to needle access disconnection detection.\nIt is especially important in instances in which a catheter enabled therapy is performed while the patient is sleeping or who otherwise may move in a manner likely to compromise catheter position and/or catheterization integrity to have an access disconnection detection system. Traditional means for attaching a catheter to a patient include the use of adhesive catheter securing devices, such as a Grip-Lok™ Universal Catheter Securement manufactured by Zefon International Inc. While a mechanical securement is beneficial, the device provides no recourse if the needle becomes dislodged despite the securement device.\nDevices that sense when a needle dislodgement occurs are also available, such as a RedSense™ device manufactured by RedSense Medical, which alarms when an absorbent pad covering the catheterization site becomes saturated with blood flowing from a wound site upon inadvertent catheter withdrawal. Disadvantages of this device include its requirement of a relatively large amount of blood to activate the device, and the device's passive audible alarm, which may not be heard by patients with hearing impairment or subject to other environmental sound source. Also, the alarm is not fed back to the dialysis machine.\nOther prior art access disconnection systems rely on current flowing through the patient's blood to indicate that a proper needle access exists. A drawback of these systems is that current must flow in many instances through high impedance or high resistance components, such as, a pump (membrane or peristaltic), a bubble trap, a valve or filter. These systems accordingly can succumb to noise, ground loops or loss of signal.\nA more robust electrical access disconnection system, in particular for blood cleansing renal therapy treatments, such as hemodialysis, hemofiltration and hemodiafiltration, in which loss of a venous needle can lead to severe blood loss, is needed accordingly."}
{"text": "This invention relates to the gun locking and security field and more particularly to a safe and secure case for a loaded hand gun that allows an authorized user rapid, unobstructed and effortless access to the hand gun.\nCrime and violence is not confined to the urban environment, but also persists throughout suburban and rural communities. As a result, there exist a large and expanding segment of the population who choose to own hand guns for self-protection, protection of others and protection of property. For this population, immediate and unrestricted access to the hand gun is crucial, since the need typically arises only in an emergency situation, such as when an intruder enters ones dwelling. Thus, anything less than free, open, unobstructed and immediate access to the hand gun defeats the primary purpose of owning the hand gun, i.e., protection.\nWide spread gun ownership has increased the number of accidents or incidents involving the use of the hand gun by an unauthorized person. Just as the population of gun owners is growing, a strong movement to protect innocent third parties, such as children, from the dangers of an accessible loaded hand gun is gaining momentum. Children in particular are intrigued by guns and have fallen victim to accidental shootings because the gun owner has taken less than adequate measures to secure the hand gun. Just as tragic is the unfortunate trend involving school shootings, which may have been prevented had the gun owner taken measures to prevent access to the hand gun by unauthorized persons. Thus, the safety-conscious segment of the population demand virtually infallible security and safety measures to prevent the unfortunate consequences that can arise from unauthorized access to a hand gun.\nUntil now, no one has been able to solve the problem of effectively marrying the two competing interests of accessibilIty and safety demanded by the seemingly opposing population segments discussed above. The present invention solves this problem by offering a hand gun case that is safe and secure from unauthorized access without sacrificing rapid and unobstructed accessibility when the need arises.\nThe present invention solves the aforementioned problems and meets the aforementioned needs by providing a secure rapidly accessible two-part impregnable gun case that uses feature recognition to allow an authorized user immediate and unobstructed access to a hand gun.\nIn the preferred embodiment, a gun case is made out of a lightweight material, such as plastic, reinforced with embedded belts that cannot be readily cut. One end of the case has an access cover that is hingably attached to the case. The case attaches Lo a fixed energized base through the use of a locking mechanism that prevents the case from being moved without unlocking it from the base. The two-part case allows the user to remove the case from the base, transport and attach the case to a fixed base in a different location, such as in a motor home or vacation property. This feature facilitates portability of the case without sacrificing safety and security.\nEnclosed in the case is a slide mechanism that supports the hand gun and projects it out of the case when the access door opens. The access cover opens when the lock is activated by a feature recognition system\"\"s identification of an authorized user\"\"s feature. It is the preferred embodiment of this invention to use fingerprint recognition, but other forms of feature recognition can be used such as voice activation.\nThe invention involves security measures used to meet the needs of the safety-conscious population. The case locks to a base that is securely attached to a surface in a location chosen by the authorized user. This prevents removal without unlocking the case. This is an improvement over the prior art, which generally cover secure cases that are not fixed to a surface. If an unauthorized person attempts to cut through the case or forcibly remove the case from the base, an audible alarm will sound. A further safety measure includes strategic placement within the case of a material, such as LEXAN, to absorb a bullet if the gun accidentally fires.\nThe invention also involves features that meet the gun owner population\"\"s demands of immediate and unobstructed access to the hand gun by using a fully enclosed spring loaded adjustable slide mechanism, which projects the gun out of the case when the access door is opened. This allows the authorized user clean, immediate and properly oriented access to the hand gun in an emergency. The gun attaches to the slide by the use of an interchangeable dowel sized to fit a particular caliber of gun. The position of the slide is adjustable depending on the size of the gun. It is desirable that the butt of the gun be located just behind the access cover so the gun projects out to maximum reach. The springs in the slide are enclosed within the slide unit and there is no support to the butt of the gun, thus once projected, the user may cleanly grab the gun.\nThe present invention also caters to both the gun owner and safety conscious population by providing immediate, unobstructed, yet secure access to the hand gun through the use of feature recognition to activate the lock securing the case access cover. In the preferred embodiment, when an emergency arises, the authorized user can touch a reader board mounted on the case. If the person\"\"s print is recognized as an authorized print, it Will activate a solenoid and unlock the access cover. This is an improvement over the prior art that teaches using a key or combination lock to prevent access by an unauthorized user. In a panic situation, fumbling with a key, possibly in the dark, or trying to remember a combination, code or sequence of button depressions is cumbersome, time consuming, and potentially life threatening. This invention also is an improvement over the prior art that teaches using a fingerprint recognition system because the prior art, requires the user to, once he or she has properly deactivated the locking mechanism, affirmatively reach into the container to retrieve the hand gun. This too, is cumbersome and time consuming depending on the orientation of the hand gun, and may also be life threatening. The feature recognition system is energized by a constant power source provided to the base, or in the event of an interruption of this source, i.e. power failure, the case contains a rechargeable battery source. A key activated manual override is provided in case there may be total failure or the energy source.\nThe present invention is an improvement over known prior art with particular reference to U.S. Pat. No. 1,557,339 issued to Sander. Sander describes a spring-loaded plate that attaches to both the butt and barrel of the gun, and projects the hand gun when the access door is opened. In Sander, the projecting springs are exposed and encumber access to the gun. Also, the gun butt is cradled by a two-piece mounting bracket that, due to the upturned ends, obstructs clean access to the gun and sacrifices speed.\nAccordingly, it is a principal object of the present invention to provide an improved gun case that meets the competing societal demands of providing a secure, tamper. resistant, safe gun case without sacrificing unobstructed access to the hand gun.\nIt is a further object of the preferred embodiment to accommodate portability of the case without sacrificing security or safety.\nThe foregoing and other objects, features and advantages of the invention will be more readily understood upon consideration of the following detailed description of a preferred embodiment of the invention, taken in conjunction with the following drawings."}
{"text": "The present invention relates to a plasma method for driving a plasma display panel. More particularly, the present invention relates to a method for driving an AC-driven plasma display panel (referred to as an AC-driven PDP hereinafter) that has a three-electrode structure and performs memory display.\nFIG. 1 is a diagram that shows the general structure of an AC-driven PDP.\nA PDP 1 comprises a pair of substrates arranged opposite to each other and a discharge gas sealed therebetween. On one of the substrates, sustain electrodes (X1 to X3) and scan electrodes (Y1 to Y3) arranged in parallel to each other are provided, and on the other substrate, address electrodes (A1 to A4) arranged in the direction perpendicular to the sustain electrodes and the scan electrodes and partitions 2 arranged in parallel to the address electrodes to define a discharge space are provided. Although only the three sustain electrodes, three scan electrodes and four address electrodes are shown in FIG. 1 for simplicity, many electrodes are actually used according to the resolution of the PDP 1.\nA display line L is formed between a sustain electrode and a scan electrode adjacent to each other. In the example in FIG. 1, the X1 electrode and the Y1 electrode form a display line L1, the X2 electrode and the Y2 electrode form a display line L2, and the X3 electrode and the Y3 electrode form a display line 3. On the other hand, adjacent pairs of the sustain electrode and the scan electrode that form the display lines form a non-display line therebetween. In the example in FIG. 1, the non-display lines are formed between the Y1 electrode and the X2 electrode, and between the Y2 electrode and the X3 electrode. In order to prevent an erroneous discharge from occurring in the adjacent display lines L1 to L3, the interval between neighboring electrodes that form a non-display line is made wider than the interval between neighboring electrodes that form a display line. Moreover, a discharge cell is formed in an area defined by a pair of the neighboring sustain electrode and the scan electrode and the address electrode that is perpendicular thereto, and a phosphor is provided in the discharge cell in order to obtain visible light.\nFIG. 2 is a block diagram that shows the general structure of the PDP apparatus shown in FIG. 5.\nThe PDP apparatus in FIG. 2 comprises the PDP 1, a data (address) driver 22, a sustain driver 23, a first (odd-numbered) scan driver 24a, a second (even-numbered) scan driver 24b, a scan pulse generation circuit 25 and an interface circuit 26 (for example, refer to Japanese Unexamined Patent Publication (Kokai) No. 10-39834).\nThe display data and the control signal from the outside of the apparatus are converted properly in the interface circuit 26 and supplied to the data (address) driver 22, the sustain driver 23, the first (odd-numbered) scan driver 24a and the second (even-numbered) scan driver 24b. The scan pulse and the sustain pulse to be applied to the scan electrode Y are generated in the scan pulse generation circuit 25 and their timings are controlled in the first (odd-numbered) scan driver 24a and the second (even-numbered) scan driver 24b by the signal from the interface circuit 26.\nSimilarly, the sustain pulse and the erasure pulse to be applied to the sustain electrode X are generated in the sustain driver 23 while being controlled by the interface circuit 26.\nA description about a method for driving the above-mentioned AC-driven PDP is given below with reference to drawings.\nA gradated display in the PDP 1 shown in FIG. 1 is performed by using the subfield driving method in which a frame is divided into a plurality of subfields and driven.\nFIG. 3 is a diagram that shows the structure of a field in the PDP shown in FIG. 1. FIG. 3 shows an example of the subfield driving method in which a field is divided into eight subfields SF1 to SF8 for a gradated display. The luminance of each field is weighted by two to the n-th power and it is possible to perform a gradated display of any level by combining proper subfields.\nIn this method, each subfield is divided into a write period (address period), a sustain discharge period and an erasure period (reset period).\nFIG. 4 is a diagram that shows the waveforms that illustrate the method for driving the PDP shown in FIG. 1. FIG. 4 shows the waveforms at the address electrode, the sustain electrodes X1 to Xn and the scan electrodes Y1 to Yn in an arbitrary subfield in a field, and each subfield is composed of the write period (address period), the sustain discharge period and the erasure period (reset period). When the scan electrode is scanned and display data is written during the write discharge period (address period), subsequent scanning is performed not to the next but the following scan electrode so that a write discharge (address discharge) is prevented from being caused to occur in the adjacent pixel successively with respect to time, and all the discharges are maintained at a time in the sustain discharge period for a light emitted display (for example, refer to Japanese Unexamined Patent Publication (Kokai) No. 2001-13915).\nAs shown in FIG. 4, in the write action (addressing) in the first half of the write period (address period), a voltage Vx is applied to the odd-numbered X electrodes X1, X3, . . . , a voltage 0V is applied to the even-numbered X electrodes X2, X4, . . . , and a scan pulse voltage −Vsc is applied to the odd-numbered Y electrodes Y1, Y3, . . . At this time, the voltage 0V is applied to the even-numbered Y electrodes. In concurrence with this, an address pulse having a voltage Va is applied selectively to the address electrode and a first discharge is caused to occur between the address electrode and the Y electrode in the selected cell in the odd-numbered display lines (L1, L3, . . . ) to be lit. Then with this discharge serving as a priming, a second discharge is immediately caused to occur between the X electrode and the Y electrode. “Address discharge” is a general term for the first discharge and the second discharge. Due to this, wall charges that enable a sustain discharge to occur are accumulated on the X electrode and the Y electrode in the selected cell in the odd-numbered display lines. When the above action is performed as far as the last odd-numbered Y electrode (Yn−1), the writing (addressing) of the selected cells in the odd-numbered display lines is completed in the first half of the write address (address period).\nNext, in the second half of the write period (address period), the voltage Vx is applied to the even-numbered X electrodes X2, X4, . . . , the voltage 0V is applied to the odd-numbered X electrodes Y1, Y3, . . . , and the scan pulse voltage −Vsc is applied to the even-numbered Y electrodes Y2, Y4, . . . , sequentially. In this way, the writing (addressing) of the selected cells in the even-numbered display lines is completed. As described above, the writing (addressing) of the selected cells in all of the display lines is completed in the first half and the second half of the write period (address period).\nIn the next sustain discharge period, a sustain pulse having a (alternating) voltage Vs is applied alternately to the Y electrode and the X electrode, a sustain discharge is caused to occur (only in the selected cells in the display lines in which the address discharge has been formed) according to the wall charges written (addressed) during the write period, as described above, and the image of a subfield in a field is displayed.\nIn the erasure period (reset period), an erasure pulse voltage VB is applied to all the sustain electrodes (X1 to Xn) to cause an erasure discharge to occur and the wall charges in the (lit) cells in the display lines, in which the sustain discharge has been caused to occur in the previous sustain period, are reduced or erased.\nHowever, in the driving method described above, the address discharge is weak in the skipped display lines (even-numbered lines in this case). As a result, a problem occurs that the light emitted display in the display line flickers or the lines appear dim.\nConcerning this problem, a description is given below with reference to FIG. 4 and FIG. 5.\nFIG. 5 is a diagram that shows how an address discharge is caused to occur when the driving method described in FIG. 4 is applied to the PDP shown in FIG. 1. For simplicity, only four sustain discharge electrodes and four scan electrodes are shown and the X1 electrode and the Y1 electrode form the display line L1, the X2 electrode and the Y2 electrode form the display line L2, the X3 electrode and the Y3 electrode form the display line L3, and the X4 electrode and the Y4 electrode form the display line L4 as shown schematically.\nAs described above, in the first half of the write period (address period), the first discharge is caused to occur between the address electrode and the Y electrode (Y1 and Y3) in the selected cell in the odd-numbered display lines (L1 and L3) to be lit, and with the first discharge serving as a priming, the second discharge is immediately caused to occur between the scan electrode Y and the sustain electrode X (between the X1 and the Y1 electrodes, and between the X3 and the Y3 electrodes).\nHowever, as the above-mentioned address discharge propagates while extending along the address electrode, it may happen that an erroneous discharge (referred to as a first erroneous discharge hereinafter) is caused to occur in the X electrode (X2) in the adjacent display line L2 adjacent to the Y1 electrode during the period of the address discharge in the odd-numbered line L1 in the first half of the write period (address period), as shown by the dotted line in the figure.\nAs a result, the address discharge during the scanning of the even-numbered line L2 in the second half of the write period (address period), which follows the scanning of the odd-numbered line L1 in the first half of the write period (address period), becomes weak and unstable, therefore, a problem occurs that the light emitted display of the display line (the even-numbered line L2) in the subsequent sustain discharge period flickers or the line appears dim.\nThe reason may be that the wall charges, which tend to decrease the potential of the sustain electrode X2 with respect to the scan electrode Y2 and the address electrode, are formed on the sustain electrode X2 due to the erroneous discharge (the first erroneous discharge), the voltage between the scan electrode Y2 and the sustain electrode X2 in the even-numbered line L2 is reduced, and the address discharge during the scanning of the even-numbered line L2 becomes weaker than the address discharge during the scanning of the odd-numbered line L1.\nIn Japanese Unexamined Patent Publication (Kokai) No. 2001-13915, a driving method has been proposed, as an improved method for driving an AC-driven PDP, in which the above-mentioned problems have been solved, the object of which is to stabilize the address discharge in the second half of the write period (address period) by increasing the voltage to be applied to the sustain electrode in the second half of the write period (address period) to recover the internal voltage that has been lowered due to the excessive wall charges caused to form by the erroneous discharge in the first half of the write period (address period), using a driving method in which either one of the odd-numbered lines and the even-numbered lines are scanned in the first half of the write period (address period) and the others are scanned in the second half of the write period (address period).\nFIG. 6 is a diagram that shows the waveforms illustrating the method for driving the PDP shown in FIG. 1. FIG. 6 shows the driving method disclosed in Japanese Unexamined Patent Publication (Kokai) No. 2001-13915, described above, wherein a voltage Vy to be applied to the sustain electrodes (X2, X4, . . . ) during the scanning of the even-lines in the second half of the write period (address period) is set to a value larger than the voltage Vx (Vx<Vy) to be applied to the sustain electrodes (X1, X3, . . . ) during the scanning of the odd-numbered lines in the first half of the write period (address period).\nAs described above, when the scanning of the odd-numbered lines and the even-numbered lines is performed separately in the first half and the second half of the write period (address period), it is possible to compensate for the amount of decrease in the potential of the sustain electrodes (X2, X4, . . . ) due to the wall charges formed on the sustain electrodes (X2, X4, . . . ) in the even-numbered lines by the erroneous discharge of the address discharge during the scanning of the odd-numbered lines in the first half of the write period (address period) by increasing the potential of the sustain electrode X with respect to the scan electrode Y and the address electrode. In this way, the address discharge can stably be caused to occur during the scanning of the even-numbered lines.\nFIG. 7 is a diagram that shows how the address discharge is caused to occur when the driving method described in FIG. 6 is applied to the PDP shown in FIG. 1.\nAs described above, the voltage Vy to be applied to the X2 electrode during the scanning of the even-numbered line L2 in the second half of the write period (address period) becomes larger than the voltage Vx to be applied to the X1 and X3 electrodes during the scanning of the odd-numbered lines L1 and L3 in the first half of the address period (Vx<Vy), therefore, the potential difference Vy+Vs between the X2 and Y2 electrodes in the even-numbered line L2 becomes larger than the potential difference Vx+Vs between the X1 and Y1 electrodes in the odd-numbered display line L1 and between the X3 and Y3 electrodes in the odd-numbered display line L3 (Vx+Vs<Vy+Vs). Due to this, the scale of the address discharge in the even-numbered line L2 becomes greater than that in the odd-numbered line L1 (by the amount of the potential difference Vy−Vx>0).\nAs a result, when the sustain discharge is caused to occur in the subsequent sustain discharge period, the scale of the sustain discharge in the even-numbered line L2 is increased by the alternating sustain pulse voltage Vs to be applied to cause the sustain discharge to occur in the selected cells in all of the display lines, and the discharge propagates to the odd-numbered line L3 adjacent to the Y2 electrode in the even-numbered line L2, as shown in FIG. 7 (a second erroneous discharge is caused to occur). This is because the wall charges, which tend to decrease the potential of the sustain electrode X3 with respect to the scan electrode Y3 and the address electrode, are formed on the sustain electrode X3 in the odd-numbered line L3 due to the second erroneous discharge, as is the same in the first erroneous discharge described above, the voltage between the scan electrode Y3 and the sustain electrode X3 in the even-numbered line L3 is reduced, and the sustain discharge in the even-numbered line L3 becomes weak and unstable. As a result, a problem occurs that the light emitted display in the display line flickers or the line appears dim."}
{"text": "1. Field of the Invention\nThe present invention relates to a photovoltaic device for converting the energy of solar light or the like into electric energy.\n2. Description of the Prior Art\nAs a photovoltaic device for converting the energy of solar light or the like into electric energy, photovoltaic devices having a variety of structures are known. One example of the photovoltaic devices is a photovoltaic device in which a transparent conductive film 2 composed of SnO.sub.2, ITO or the like is provided on a transparent substrate 1, a photoelectric conversion layer 3 in which a p-type (or n-type) semiconductor layer 3a (or 3c), an intrinsic (i-type) semiconductor layer 3b, and an e-type (or p-type) semiconductor layer 3c (or 3a) are laminated is provided on the transparent conductive film 2, and a back electrode 4 composed of a high-reflective metal such as Ag or Al is provided on the photoelectric conversion layer 3.\nIn such a photovoltaic device, light is introduced into the above described photoelectric conversion layer 3 through the transparent conductive film 2 from the side of the above described transparent substrate 1 whereby carriers are generated in the above described intrinsic layer 3b to carry out a photoelectric conversion.\nFurthermore, there has been developed in recent years a photovoltaic device so constructed that the interface between the transparent conductive film 2 and a photoelectric conversion layer 3 is so textured as to have an irregular shape so that reflection of light from the interface between the transparent conductive film 2 and the photoelectric conversion layer 3 is reduced when the light is introduced into the photoelectric conversion layer 3 through the transparent conductive film 2, thereby increasing the amount of the light introduced into the photoelectric conversion layer 3, and the light introduced into the photoelectric conversion layer 3 is scattered, thereby increasing the optical path length of the light passing through the photoelectric conversion layer 3, whereby many carriers are generated in the intrinsic layer 3b, thereby increasing the conversion efficiency as shown in FIG. 2.\nIf the interface between the transparent conductive film 2 and the photoelectric conversion layer 3 is formed in an irregular shape, however, the internal electric field in the photoelectric conversion layer 3 comprising a semiconductor junction which constituted by the p-type semiconductor layer 3a, the intrinsic layer 3b, and the n-type semiconductor layer 3c becomes non-uniform in the surface direction of the layer. As shown in FIG. 3, in a portion of the photoelectric converting layer 3 which is in contact with a recess 2a of the transparent conductive film 2 where much light is absorbed, the density of lines of electric force indicated by broken lines is decreased toward the transparent conductive film 2. Consequently, the interval between equipotential lines indicated by dot and dash lines is gradually increased toward the transparent conductive film 2, so that the internal electric field in the photoelectric conversion layer 3 in the portion which is in contact with the recess 2a of the transparent conductive film 2 becomes weak. In the portion of the photoelectric conversion layer 3 where the internal electric field becomes weak, many carriers recombine, whereby fill factor of the photovoltaic device is decreased, thereby preventing an improvement in conversion efficiency."}
{"text": "The present application relates to an optical element, a replica substrate configured to form an optical element, and methods for producing the optical element and the replica substrate. Specifically, the present application relates to an optical element having many structures each being in the form of a projection or depression on a surface of a base.\nIn an optical element including a transparent substrate composed of a material such as glass or plastic, in order to reduce surface reflection of light to improve light transmission characteristics, a method for forming fine and dense projections and depressions (subwavelength structure: moth-eye structure) on a surface of the optical element is employed. Generally, in the case where periodic projections and depressions are provided on a surface of an optical element, light passes through the projections and depressions and is diffracted, significantly reducing the amount of light components that travel in straight lines. However, in the case where the pitch of the projections or depressions is smaller than the wavelength of light transmitted, diffraction does not occur. For example, in the case where each of the projections is in the form of a cone, an effective antireflection effect and excellent transmission characteristics are obtained for single-wavelength light corresponding to the pitch and the height of the projections.\nFor example, Non-Patent Document 1 (NTT Advanced Technology Corporation site (accessed May 21, 2006), “Master Mold for Antireflective (Moth Eye) Structure Independent of Wavelength”) describes a method for producing the foregoing optical element. A photoresist on a Si substrate is formed into an uneven photoresist pattern by electron beam lithography. The Si substrate is subjected to etching with the resulting uneven photoresist pattern as a mask, thereby forming a Si master mold having subwavelength structures each being in the form of a fine cone (pitch: about 300 nm, depth: about 400 nm) as shown in FIG. 15.\nThe resulting Si master mold described above can also have an antireflection effect on light having a wide wavelength range. As shown in FIG. 16, the arrangement of the subwavelength structures in a hexagonal lattice results in an extremely high antireflection effect (reflectivity: 1% or less) in the visible region (see FIG. 17). In FIG. 17, 11 represents the reflectivity of a flat portion of the Si master mold, and 12 represents the reflectivity of the patterned portion of the Si master mold.\nAs shown in FIG. 18, a Ni electroformed stamper of the Si master mold is produced. As shown in FIG. 19, the stamper has an uneven pattern in a predetermined region RI of a surface of the stamper, the pattern being the inverse of the pattern of the surface of the Si master mold. The uneven pattern is transferred to a transparent polycarbonate resin with the stamper, thereby affording a target optical element (replica substrate). The resulting optical element also has a high antireflection effect (reflectivity: 1% or less) (see FIG. 20). In FIG. 20, 13 represents the reflectivity of a portion other than the pattern, and 14 represents the reflectivity of the pattern. However, Non-Patent Document 1 does not describe the transmission characteristics of the optical element or improvement in the light-guiding performance of a display or the like with the optical element.\nIn recent years, attempts have been made to apply uneven surfaces of subwavelength structures to antireflection elements and high-transmittance elements used for various optical devices such as displays, optoelectronics devices, optical communication devices, solar cells, and lighting systems. Hereinafter, liquid crystal displays as an example of the application of the uneven surfaces of subwavelength structures will be described below.\nIn recent years, ultrathin displays, such as liquid crystal displays (LCDs) and plasma display panels (PDPs), have been practically used instead of cathode-ray tube (CRTs) displays. In particular, LCDs are spreading at an accelerating pace because of low-power-consumption operation and a reduction in the cost of large color LCD panels.\nIn liquid crystal displays, backlight systems in which transmissive color LCD panels are illuminated from the back side of the panels by backlights to display color images are mainly employed. As light sources of backlight systems, cold-cathode fluorescent lamps (CCFLs) that emit white light are often used.\nCold-cathode fluorescent lamps contain mercury in their fluorescent tubes and may adversely affect the environment. Thus, as light sources of backlight systems, light sources to replace cold-cathode fluorescent lamps may be needed. Since blue-light-emitting diodes have been recently developed, light-emitting diodes (LEDs) emitting red light, green light, and blue light, which are the primary colors of light, have been commercially available. Thus, LEDs hold promise for light sources to replace cold-cathode fluorescent lamp. Mixing red light, green light, and blue light emitted from LEDs results in white light with excellent color purity.\nThe use of LEDs as light sources of backlight systems enhances the color purity of colored light through color liquid crystal display panel. Thus, the color reproduction range is expected to extend to the extent specified by the NTSC (national television system committee) system or to the extent more than specified.\nIn the case where LEDs emitting red light, green light, and blue light are used as light sources of backlight systems configured to illuminate color LCD panels, a technique for efficiently transmitting each colored light may be required. As described above, therefore, attempts have been made to apply uneven surfaces of subwavelength structures to LCDs.\nFor example, Japanese Patent No. 3723843 discloses an antireflection filter (subwavelength structure) formed by subjecting a resist film to electron beam lithography and then development to form a resist pattern and directly subjecting a second semiconductor layer to etching with the resist pattern using a predetermined etching gas. The patent document also describes that a semiconductor light-emitting element having the antireflection filter has an emission intensity 30% higher than that of a semiconductor light-emitting element not having the antireflection filter.\nHowever, electron beam exposure disadvantageously requires a long operation time and thus is not suitable for industrial production. For example, in the case where an electron beam with a beam current of 100 pA, which is used for the finest pattern, is used for a resist such as calixarene requiring a dose of several tens of mC/cm2, the exposure of a square with a side length of 200 μm is not completed even after 24-hour exposure. In the case where the area of a 2.5-inch display (50.8 mm×38.1 mm), which is commonly used for mobile phones, is exposed, it takes about 20 days.\nA method described in Non-Patent Document 2 (National Institute of Advanced Industrial Science and Technology site. http://aist.go.jp/aist_i/press_release/pr2006/pr20060306/pr20060306.html (accessed May 21, 2006), “Development of Desktop Apparatus Enabling Nanometer-Scale Microfabrication”) utilizes optical-disk-recording technology. Thus, a large-area optical element can be produced at high speed and low cost. However, the resulting optical element produced by the method has poor wavelength-dependent characteristics of reflectivity and a reflectivity exceeding 1%; hence, the optical element is not practical for an antireflection structure. This may because the density (aperture ratio) of the nanodot pattern is low (50% or less).\nIn addition to the foregoing subwavelength structures, attempts have been made to apply inorganic multilayer films to LCDs. For example, Japanese Unexamined Patent Application Publication No. 2006-145885 discloses that low-refractive-index dielectric layers composed of silicon oxide (SiO2) and high-refractive-index dielectric layers composed of niobium pentoxide (Nb2O5) are stacked to form a 24-layer optical filter having the foregoing optical properties, thereby increasing transmittances of red light (640 nm), green light (530 nm), and blue light (450 nm) to 80%, 80%, and 50%, respectively. However, such an optical filter does not sufficiently improve transmittance characteristics. Thus, further improvement in transmittance characteristics, in particular, improvement in transmittance characteristics of blue light is desired."}
{"text": "1. Field of the Invention\nThe present invention is directed to an x-ray tube of the type having an evacuated housing rotatable around a rotational axis with a drive arrangement in which, firmly connected to the housing, a cathode that emits electrons and an anode onto which the electron beam accelerated with an electrical field, strikes are arranged, and having a non-rotating (stationary) electromagnetic system for the deflection and focusing of the electron beam that includes a number of current-permeated coil elements.\n2. Description of the Prior Art\nAn x-ray tube of the above type is described in U.S. Pat. No. 4,993,055. In this rotating bulb tube, the elements required for generating the x-radiation, i.e. the cathode and the anode, are firmly connected to the housing and rotate together with it during the operation of the x-ray tube. In order to be able to guide the beam on a predetermined path during the rotation so that the path proceeds quasi stationarily and always strikes the rotating anode in the same focal spot, a stationary electromagnetic system is provided that is fashioned for the deflection as well as for the focusing of the electron beam. According to this patent, this electromagnetic system has four coils, known as quadrupole coils, that are arranged around the housing of the tube and that enable operation in a combined deflection/focusing mode on the basis of corresponding control. In this way, it is possible to control the electron beam and the generated focal spot on the anode from the exterior within broad ranges.\nThe arrangement and fastening of the individual coil elements represents a significant problem, since the coil elements are to be provided in the region of the housing in as simple a way as possible, but a high alignment precision must be achieved at the same time. U.S. Pat. No. 4,993,055, wherein rectangular coils having large dimensions are employed, does not disclose a satisfactory solution."}
{"text": "Vinyl chloride polymers and copolymers are known to deteriorate under the influence of heat, so resulting in alteration of the physical properties thereof and discoloration which precludes use in transparent and lightly colored articles. This thermal, nonoxidative degradation involves the sequential loss of hydrogen chloride from adjacent monomer units and results in the formation of conjugated polyene structures which impart undesirable color to the vinyl chloride polymer.\nEfforts to obviate these limitations have included the use of stabilizing compositions of the MY.sub.2 type wherein M is a metal cation, for example R.sub.2 Sn.sup.+2 (R=alkyl), Ba.sup.+2, Cd.sup.+2, Zn.sup.+2, Pb.sup.+2, Ca.sup.+2, and so forth, and Y is an organic anion. The most effective stabilizers arising from this class of compounds are those classified as organotin compounds. It has been theorized that these MY.sub.2 stabilizers react with hydrogen chloride, and since hydrogen chloride is known to catalyze PVC degradation, the efficacy of these stabilizers is due in part to hydrogen chloride scavenging. However, the stabilizers are also believed to have the ability to undergo rapid and selective metathetical reactions which destroy certain PVC structural anomalies that are the active degradation sites. Accordingly, in addition to reducing the rate of color development in PVC by interrupting polyene growth, the MY.sub.2 stabilizers should manifest a true chemical stabilization effect by reducing the rate of formation of total chloride.\nDespite conflicting theories advanced by workers in the art relative to the described stabilization mechanism, the MY.sub.2 stabilizers have been used extensively in industry. Although considered satisfactory from the standpoint of stabilization, toxicity of the heavy metal residues and ecological considerations have stimulated further evaluation of the subject compositions and have generated a search for alternatives.\nIn accordance with the present invention, this end is attained by means of a novel process which yields a product that is free from toxic heavy metal residues characteristic of most of the prior art. Additionally, the polymer so prepared evidences enhanced stability with respect to the unstabilized compositions, and such stability has been found comparable to that provided by many prior art vinyl chloride stabilizers.\nBriefly, the inventive technique involves heating a vinyl chloride polymer in solution with an organotin compound, precipitating the vinyl chloride polymer from the reaction product, and subjecting the resultant precipitate to solvent extraction to remove heavy metal residues."}
{"text": "The invention relates to an openable twin tank toilet flush box of multi-outlets using recycled water. Openable tanks will make it easy to clean, maintain and stop residues in the flush box under hand, washing tubs by using recycled water from hand washing tubs and bathtubs to be used in toilet flush boxes of various types."}
{"text": "Due to demand for miniaturized, highly functional electronic products in the semiconductor industry, a variety of techniques for providing high-capacity semiconductor packages are being researched and developed.\nIn general, since wafer level packages are packaged on a wafer level, packaging costs can be considerably reduced according to the number of semiconductor die produced for each wafer. However, when a stress buffer is removed (e.g., a substrate serving as a stress buffer), the wafer level packages have poor reliability compared to packages with the stress buffer in place.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with the present disclosure as set forth in the remainder of the present application with reference to the drawings."}
{"text": "The thermoresistive materials used for measuring temperature in the field of application up to 1400° C. are mechanically unstable and, thus, typically cannot be used as self-supporting temperature sensor elements. Therefore, they are typically mounted in protective tubes or on or between substrates. These are usually ceramic substrates. Known temperature sensors that are suitable for exhaust gas are thermoelements, which are typically made of precious metal wires of Pb/PtRh or Ni/CrNi, whose connecting legs are mutually insulated in ceramic tubes and whose contact point is protected by a metal casing or ceramic casing or by being welded into the metal casing of a sheathed type glow plug. Further known are temperature sensors designed as thick or thin layer elements, where the temperature-sensitive functional layer is vapor deposited or sintered on or between the substrates. This requires a certain (thermal) inertia of the thermoelements on the part of the respective carrier material.\nFurthermore, it is quite difficult to measure temperature during the combustion processes in the combustion chamber of combustion engines. Particularly in the case of modern four-valve direct injection engines, there is often no possibility to provide an additional borehole for the access to the combustion chamber for a temperature sensor.\nMoreover, the temperatures or the temperature ranges from −40 to 1400° C. to be measured, in conjunction with an aggressive atmosphere in the form of hot gas place extremely high demands on corresponding temperature sensors."}
{"text": "1. Field of the Invention\nThe present invention relates to a boring and conduit/pipe system and more particularly pertains to forming a hole through a wall and for pulling a linear member through a hole.\n2. Description of the Prior Art\nThe use of boring and conduit/pipe systems of known designs and configurations is known in the prior art. More specifically, boring and conduit/pipe systems of known designs and configurations previously devised and utilized for the purpose of boring through known methods and apparatuses are known to consist basically of familiar, expected, and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which has been developed for the fulfillment of countless objectives and requirements.\nBy way of example, U.S. Pat. No. 5,888,035 issued Mar. 30, 1999 to Cutler relates to a system for installing a length of cable within a wall. U.S. Pat. No. 6,561,269 issued May 13, 2003 to Brown relates to a canister, sealing method and composition for sealing a bore hole. U.S. Pat. No. 5,242,026 issued Sep. 7, 1993 to Deken relates to a method and apparatus for drilling a horizontal controlled borehole in the earth. U.S. Pat. No. 5,133,418 issued Jul. 28, 1992 to Gibson relates to a directional drilling system with eccentric mounted motor and biaxial sensor and method. U.S. Pat. No. 4,953,638 issued Sep. 4, 1990 to Dunn relates to a method and apparatus for drilling a horizontal controlled borehole in the earth. U.S. Pat. No. 4,942,930 issued Jul. 24, 1990 to Millsapps relates to a lubrication system for an earth boring drill bit and method for filling and retrofit installing thereof. U.S. Pat. No. 4,143,723 issued Mar. 13, 1979 to Schmotzer relates to a carbide tipped drill bit for boring holes in concrete and steel. Lastly, U.S. Pat. No. 3,737,245 issued Mar. 30, 1973 to Mater relates to a drill bit for pole boring machine.\nWhile these devices fulfill their respective, particular objectives and requirements, the aforementioned patents do not describe boring and conduit/pipe system that allows forming a hole through a wall and for pulling a linear member through a hole.\nIn this respect, the boring and conduit/pipe system according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in doing so provides an apparatus primarily developed for the purpose of forming a hole through a wall and for pulling a linear member through a hole.\nTherefore, it can be appreciated that there exists a continuing need for a new and improved boring and conduit/pipe system which can be used for forming a hole through a wall and for pulling a linear member through a hole. In this regard, the present invention substantially fulfills this need."}
{"text": "Many aircraft in use today require the presence of an auxiliary power unit at one stage or another during their operation. Generally speaking, the function of the auxiliary power unit is to provide electrical and/or hydraulic power to aircraft systems when the aircraft's main engines are not in operation.\nAdditionally, in aircraft operating environments, it is also desirable in many instances to have a source of compressed air that may be utilized, for example, in starting the aircraft's main propulsion engines or in providing compressed air to aircraft systems that ordinarily would receive compressed air as bleed air from their main propulsion engines.\nFurther, not infrequently, such environments also require an air turbine starter for initiating operation of a turbine and/or an air turbine starter motor which is much the same as an air turbine starter but is adapted to be operated more or less continuously.\nNeedless to say, to provide all the equipment necessary to perform the foregoing functions is quite expensive, and where such equipment is on board the aircraft, it lessens the load carrying capability of the aircraft, all other things being equal, because of the weight of such equipment.\nThe present invention is directed to overcoming one or more of the above problems."}
{"text": "The increased popularity of ice skating in cold weather areas has led to great competition for the limited resources of the public and private ice skating venues presently available. As a result, the costs of using private ice skating facilities has increased dramatically and public rinks must generally maintain long hours to accommodate the large numbers of hockey teams, figure skaters and the desires of the general public for time on the ice. Often times a hockey club or those practicing for competitive figure skating will be forced to accept practice time in the very late evenings or very early morning hours.\nTraditionally, skaters have used a frozen lake or pond on which to skate in the winter time. Where a lake or pond is not available a field has been flooded and allowed to freeze-up to provide a skating surface. This is done by flooding an area with water until the ground is saturated and the local water table rises above ground level and allowing the water to freeze. However, because of variable weather conditions in the winter months, skating on frozen ponds and lakes is extremely dangerous and the ice on flooded fields may melt causing a loss of all of the water poured onto the field.\nPresently, the method used to make an outdoor skating rink is to create some sort of perimeter, place a liner in the perimeter, fill it with water and let it freeze. The perimeters have generally been constructed by driving dimensional lumber, such as two-by-twos or two-by-fours, part way into the ground, as a structure, and then attaching plywood, using nails or screws, to the structure to form the perimeter. On flat fields such a system, while requiring a great deal of labor to construct, works well. However, where there is a slope to the ground, as in most areas, the perimeter must be made so that the walls on the downward part of the slope are high enough to contain the deepest portion of the water. Often times during construction of the rink structure, because a field will often look more level than it is, the height of the water is not taken into account until water is poured into the rink. After the water is poured into the rink it is difficult to remove a section, or entire wall, of plywood and/or to drive more or larger dimensional lumber into the ground for support.\nAlso, such construction generally must begin early in the season while the ground is still soft enough to drive dimensional lumber into the ground. Because of the variability of weather, it has often been found that the effort of constructing such a rink is lost because a winter season is not cold enough to allow skating often enough or long enough to justify the effort of such construction. As a result, owners of these rinks are discouraged from exerting the effort to construct such rinks in future years. Also, driving large pieces of lumber into the ground, to construct the structure of the rinks, causes considerable damage to a field, or back yard, which must be repaired in the summer time.\nIt is therefore an object of the present invention to provide an ice skating rink structure that is easy and quick to construct.\nIt is another object of the present invention to provide an ice skating rink structure that can be easily used on level or sloped ground.\nIt is a further object of the present invention to provide an ice skating rink structure that can easily accommodate various sizes of plywood, or other boards, and that can allow the easy removal and replacement of boards.\nIt is another object of the present invention to provide an ice skating rink structure which can be easily constructed at any time during the winter season without regard to the condition of the soil.\nIt is another object of the present invention to provide a skating rink structure that can be installed and removed without causing great damage to a field or lawn.\nOther objects and advantages of the present invention will become apparent as the description proceeds."}
{"text": "The present invention relates to an ionization smoke detector connected to an alarm circuit triggered by a sensor in the detector when a predetermined smoke level is exceeded. The detector has particular utility as a fire detector in residences since it is battery-operated and the alarm circuit provides warnings for both fires and approaching weak battery conditions.\nIonization smoke detectors such as shown in U.S. Pat. No. 3,728,706, are already well known in the art. Such detectors operate upon the principle that an ionization current catalyzed by a radiation source is affected by the presence of particulate matter such as that found in aerosols or smoke emanating from a fire. The particulates reduce the ionization current and such a reduction in current can be detected and correlated with the density of the particulate matter to provide a fairly accurate indication of a smoke condition.\nIt is desirable to operate such a detector with a battery, especially in residential homes and apartments. To provide a prominent alarm signal, a pulsating circuit energized by the battery may be triggered by the smoke detector. A pulsating circuit reserves battery strengths so that the alarm may continue for an extended period of time even after the life of the battery is partially expended.\nIn view of the reliance that may be placed upon a smoke detector, particularly in a residential environment, safety considerations must be given careful consideration. The limited life of a battery might suggest that it is not suitable for an alarm device; however, it is also recognized that fire can destroy or cut off outside electrical sources before the fire is detected. The battery-operated system functions independently and can be safely relied upon if it includes low-voltage detection circuitry for sounding the alarm when the battery level is approaching a weakened or marginal condition.\nIt is, accordingly, a general object of the present invention to disclose a battery-operated ionization smoke detector connected to an alarm circuit possessing features insuring safe operation in the residential environment."}
{"text": "The present invention relates generally to multiplexing of communications signals and more particularly to the wavelength division multiplexing of optical signals in fiber optics communications.\nWavelength Division Multiplexing (WDM) is a way of increasing the capacity of an optical fiber by simultaneously operating at more than one wavelength within a single optical fiber. Multiple optical signals of different wavelength are transmitted in the same direction over one strand of fiber, and the signals are later separated by wavelength at the distant end. In order to establish some standards for WDM operations, the International Telecommunication Union (ITU) has proposed successive systems of standardized frequencies to be used as channels for optical telecommunications, with each system incorporating more and more channels, usually with smaller and smaller wavelength separation between the channels. This system of channels is spoken of as the xe2x80x9cITU gridxe2x80x9d and currently includes 80 channels utilizing a wavelength range centered around 1550 nm (193,300 GHZ ) with a channel spacing of approximately 0.4 nm (50 GHZ). An ITU grid with 50 GHZ frequency spacing is shown in FIG. 1 with the fringe order pattern from a Fabry-Perot interferometer using an etalon of appropriate parameters superimposed on the grid. There has also been a system proposed which uses channel spacing of 25 GHZ (0.2 nm). It will be easily apparent that the smaller the spacing is between channels, the more sensitive a multiplexed channel system will be to fluctuations that will cause the channel frequencies to drift away from the target grid frequencies. In a system where the frequency spacing is 100 GHz, the variation of from the grid frequency of 0.01 nm (1.2 GHz) would be undesirable, but perhaps may not be disastrous. In a 50 GHz system, a 0.01 nm variation would impair the performance of the system, and in a 25 GHz system, such a deviation would likely be disabling.\nAn interference filter is a type of tool that is often used to separate multiple wavelengths of light that are included in a beam of light. A Fabry-Perot Interferometer is one type of interference filter that is often used for wavelength filtering and separation. Interference filters operate by providing a pair of mirrored surfaces with a cavity formed between them. Incident light undergoes multiple reflections between the mirrored surfaces, which typically reflect greater than 95% of the light at each surface. The incident and reflected wave interfere with each other constructively or destructively depending on their phase relationship. Where there is no phase difference between successive wave, constructive interference is produced and a maximum is produced in the transmission portion. Where the phase are 180 degrees out of phase, destructive interference occurs and a minimum is transmitted. A maximum occurs when the round trip optical path is some integer multiple of whole wavelengths, and also depends on the thickness of the cavity (d), the index of refraction of the cavity material (n), and the angle of incidence (0), which are related by the formula:\n2d n cosxcex8=mxcex,\nwhere m is an integer, often termed the order number and xcex is the wavelength of the light. The parallel rays of each wavelength are often focused by a lens in order to produce a familiar ring pattern. The result is a series of transmission peaks of separated wavelength. The separation distance between adjacent peaks is equally spaced when plotted with respect to inverse wavelength, and is called the Free Spectral Range (FSR).\nEtalons are special Fabry-Perot interferometers which have fixed spacing between the reflective surfaces, thus the thickness of the cavity d is therefore not subject to direct parallel variation. However, the etalon may be tilted, changing the angle of the etalon relative to the angle of incidence of the light beam, which thus increases the optical path length. This allows the etalon to be xe2x80x9ctunedxe2x80x9d over a limited range to alter the peak transmission wavelengths. FIGS. 2 and 3 show etalons as used in the prior art. FIG. 2 shows the incident light striking the etalon at an angle xcex8, and FIG. 3 shows the same etalon tuned at the same angle xcex8, thereby increasing the optical path length n to angle tune the etalon.\nHowever, angle tuning of an etalon introduces other problems caused by the insertion loss due to the variation in angle. Besides the difficulties of producing very tiny variations in angle, when the etalon is tuned at a small angle, the output beam can become oblong in shape, with non-uniform beam intensity distribution. As this angle increases, this effect becomes more pronounced. When optics are used to collect the output light from the etalon, a large insertion loss variation is often seen. This variation is typically from 1-4 dB.\nThe variation of etalon insertion loss also commonly causes the operating point of the output spectrum to shift by as much as 10 pm (0.01 nm=1.2 GHz). As discussed above, errors of this magnitude can seriously interfere with operation of systems which use 25 Ghz frequency separations and even with 50 GHz systems.\nThus, there is a great need for a etalon which is usable in a multi-channel wavelength locking system which does not produce such large variations in insertion loss, beam quality, and wavelength shift.\nAccordingly, it is an object of the present invention to provide a multi-channel wavelength locking system with better ITU setting repeatability.\nAnother object of the invention is to provide a multi-channel wavelength locking system with improved temperature stability.\nAnd another object of the invention is to provide a multi-channel wavelength locking system which produces smaller variations in insertion loss.\nA further object of the present invention is to provide a multi-channel wavelength locking system which produces smaller variations in beam quality.\nAn additional object of the present invention is to provide a multi-channel wavelength locking system which produces less wavelength shift.\nYet another object of the present invention is to provide a multi-channel wavelength locking system which can be used with wavelengths which are separated by as little as 25 GHz or less.\nA yet further object of the present invention is to provide a tunable etalon which has reduced manufacturing costs due to more relaxed dimensional tolerances in the parts.\nBriefly, one preferred embodiment of the present invention is a wavelength locker for providing wavelength-locked multi-channels signals, including a number of radiation sources providing beams of radiation at wavelengths separated by a predetermined spacing, each spaced wavelength being a channel which is connected to a multiplexer. One or more wavelength lockers produce equally-spaced spectral lines which are tunable wavelength locker includes one or more etalons which includes a gas-tunable medium having a variable optical index of refraction. The etalons may either be transmissive, or reflective with respect to which surface the major portion of the light exits. The etalons are tuned by varying the pressure or composition of the gas-tunable medium until the spectral lines of the etalon align with the ITU grid. The gas properties are then fixed by at least temporarily sealing the etalon enclosure.\nIn use in a multi-channel system, beam splitters divert a portion of the beams of radiation into the wavelength lockers. Detectors then receive the spectral line output of the wavelength lockers and detects shifts in wavelength of the spectral line output to generates control signals. The wavelength of each of the radiation sources is then adjusted in response to the control signals.\nThe output of the detector may also be normalized in order to isolate the variations which are attributable to shift in wavelength. The channels may be sequentially sampled and the outputs adjusted by a single wavelength locker or there may be separate wavelength lockers for each input channel.\nA method for locking the wavelengths of a plurality of radiation sources to a wavelength reference is also included.\nAn advantage of the present invention is that it provides improved ITU setting repeatability.\nAnother advantage of the present invention is that it provides improved temperature stability.\nAnd another advantage of the present invention is that it produces very small variations in insertion loss and beam quality.\nA further advantage of the present invention is that it can be used with wavelengths which are separated by as little as 25 GHz or less.\nA yet further advantage is that the present invention can provide a tunable etalon which has reduced manufacturing costs due to more relaxed dimensional tolerances in the parts.\nThese and other objects and advantages of the present invention will become clear to those skilled in the art in view of the description of the best presently known mode of carrying out the invention and the industrial applicability of the preferred embodiment as described herein and as illustrated in the several figures of the drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital broadcast system for transmitting/receiving digital broadcast data, and a data processing method for use in the same.\n2. Discussion of the Related Art\nA variety of digital broadcast standards have been widely used throughout the world. For example, a vestigial sideband (VSB) transmission scheme has been used as a digital broadcast standard for the North America and the Republic of Korea. This VSB scheme is based on a single-carrier scheme, such that a reception (Rx) performance of a reception (Rx) system may deteriorate under poor channel environments. Specifically, a portable or mobile broadcast receiver requires much higher resistance to a channel variation or noise. So, if mobile service data is transmitted according to the VSB transmission scheme, a reception (Rx) performance becomes deteriorated."}
{"text": "1. Field of the Invention\nThe present invention relates to rectal balloons that are used for immobilizing the region surrounding the prostate.\n2. Description of the Related Art\nTreatment of prostate cancer using radiation therapy is difficult due to the prostate's position near radiation-sensitive tissues, and is further complicated by prostate motion. Adenocarcinoma of the prostate commonly occurs in the posterior portion of the prostate gland, which is in very close proximity to the rectal wall. To date, external beam radiation treatment, urethrograms, CT scans and magnetic resonance imaging (MRI) have all been used to visually localize the prostate, as well as the normal critical structures in the surrounding area.\nU.S. Pat. No. 5,476,095, issued on Dec. 19, 1995 to Schnall et al., proposes an insertable pickup probe for use in providing diagnostic MRI images. The pickup probe, in its preferred embodiment, is proposed for use in imaging the male prostate and comprises an elongated shaft supporting an inflatable patient interface balloon at its distal end. The interface balloon comprises an inner balloon and an outer balloon, between which a receiving coil is positioned. A lumen for air supply is provided in the shaft for expanding the inner balloon against the outer balloon to place the receiving coil in close proximity to the area of interest in order to provide MRI images.\nTypically, the planning of radiation therapy for the treatment of prostate cancer involves the patient undergoing a CT-based simulation scan of the pelvis to determine the location of the prostate gland. In the simulation phase, the patient is placed on CT equipment that is preferably similar to the radiation treatment equipment (except that it does not generate the high energy radiation beam). The simulation equipment is positioned to simulate the delivery of the sequence of treatment beams prescribed by the treating oncologist. Normally, during the simulation procedure, CT images are acquired. These CT images allow the oncologist to locate the position of the tumor and help to facilitate the composition of a radiation treatment plan. This treatment plan delineates the positions of the radiation equipment components for delivery of the treatment beams.\nDuring the actual treatment phase, the patient is placed in the same position on the treatment equipment as in the simulation scans. Radiation-emitting devices are generally known and used for radiation therapy in the treatment of patients. Typically, a radiation therapy device includes a gantry, which can be swiveled around a horizontal axis of rotation in the course of a therapeutic treatment. A linear accelerator is located in the gantry for generating a high-energy radiation beam for therapy. During treatment, the radiation beam is provided by this equipment and is delivered to the patient at the precise location as delineated by the physician during simulation. A further feature of radiation therapy involves portal images, which are commonly used in radiation therapy to verify and record the patient tumor location. Portal images include manual (film) and electronic images (EPI) taken before and/or after the treatment.\nDuring external beam radiation therapy, radiation is directed to the target prostate which is near the rectal wall. A misdirected radiation beam may perforate the rectal wall causing radiation proctitus (rectal bleeding). This toxicity is related to the total radiation dose prescribed and the volume of the anterior rectal wall receiving a high radiation dose. A major factor limiting radiation oncologists' attempts to reduce the volume of the anterior rectal wall receiving a high radiation dose is the position of the prostate gland as well as the intrinsic motion up to 5 mm in the anterior to posterior direction caused by rectal peristalsis. Accordingly, oncologists generally will add a margin to the radiation field in order to ensure that the entire prostate gland receives the prescription dose. This margin is typically on the order of 5 to 15 mm. As a consequence, lower doses of radiation may need to be used so as not to overexpose radiation sensitive structures. However, this may lead to inadequate radiation treatment and a higher probability of local cancer recurrence.\nU.S. Patent Publication No. 2003/0028097, published on Feb. 6, 2003 to D'Amico et al., proposes an insertable probe for immobilizing a region of interest during staging and radiation therapy thereof. In particular, a balloon is proposed having a rectangular cross section connected to a shaft. The shaft extends to an end of the balloon so as to allow fluid flow through an interior of the shaft and into the balloon so as to selectively inflate the balloon once the balloon is installed into the rectal cavity. The balloon, shaft and handle are bonded together so that they move radially as a single unit when torque is applied. A syringe is provided which connects the shaft and serves as an air pump to deliver a volume-limited amount of air to the air lumen of the shaft to the balloon. A stop cock is provided to maintain the air within the balloon.\nOne of the problems with the subject of U.S. Patent Publication No. 2003/0028097 is the discomfort associated with installing the rectal balloon within the rectal cavity. In particular, a relatively sturdy and wide diameter shaft is connected to a relatively large thick-walled balloon. Because the balloon is not supported by anything other than by the shaft, the balloon is formed of a relatively rugged and thick material. Because of the relatively large size of the shaft and the thick material of the rectangular-cross section balloon, the installation of the rectal balloon creates a large amount of discomfort for a patient. It is often difficult for the medical personnel to know exactly how far within the rectum the balloon has been installed. It is difficult to achieve a standardized and fixed position of the balloon during each and every use. The medical personnel must generally approximate the desired position of the balloon within the rectal cavity. As such, a need has developed whereby the rectal balloon can be formed of a minimal diameter shaft and of a balloon of relatively thin material.\nWhen the rectal balloon of U.S. Patent Publication No. 2003/0028097 is in an inflated condition, the outer surface is generally round. As such, the prostate will tend to balance on the curved surface rather than be properly seated thereon. Since seating is important for proper use, this device requires that the physician approximate a seated position rather than providing any feedback of the seated position. When the balloon is in a curved inflated condition, the prostate will have a tendency to slide to one side of the balloon. As such, a need developed to provide a rectal balloon that retains the prostate in a proper seated position when the balloon is in a fully inflated condition.\nAs discussed above, a very important consideration when treating patients using radiation therapy is that the proper dose of radiation reaches the treatment site. This is very important whether the treatment method utilizes implanted radiation seeds or external beams of radiation. Excessive dosing of the patient can lead to severe side effects including impotence and urinary incontinence. A proper treatment plan should deliver an adequate amount of radiation to the treatment site while minimizing the dose delivered to the surrounding tissues.\nU.S. Pat. No. 6,963,771, issued on Nov. 8, 2005 to Scarantino et al., describes a method, system and implantable device for radiation dose verification. The method includes (a) placing at least one wireless implantable sensor in a first subject at a target location; (b) administering a first dose of radiation therapy into the first subject; (c) obtaining radiation data from the at least one wireless implantable sensor; and (d) calculating a radiation dose amount received by the first subject at the target location based on the radiation data obtained from the at least one wireless sensor during and/or after exposure to the first administered dose of radiation to determine and/or verify a dose amount of radiation delivered to the target location.\nU.S. Pat. No. 7,361,134, issued on Apr. 22, 2008 to Rozenfeld et al., proposes a method of determining the dose rate of a radiation source including locating three or more detectors in the vicinity of a source. Each of the detectors provides an output indicative of the amount of radiation received from the source and determines the location of the source from at least some of the detector outputs. International Pub. No. WO 2008/148150, published on Dec. 11, 2008 to Rozenfeld, proposes a semiconductor radiation sensor.\nU.S. Patent Publication No. 2009/0236510, published on Sep. 24, 2009 to Lacroix et al., proposes a radiation dosimeter for measuring a relative dose of a predetermined radiation type within a detection region by using a plurality of scintillating optical fibers. U.S. Patent Publication No. 2006/0094923, published on May 4, 2006 to Mate, proposes a marker comprising a wireless transponder configured to wirelessly transmit a location signal in response to a wirelessly transmitted excitation energy.\nA significant cause of patient discomfort associated with rectal balloons of the prior art is the buildup of gas pressure when the balloon is inserted into the rectum. This buildup of gas pressure can also affect the ability to properly seat the balloon in a desired position.\nThe above discussed U.S. Pat. Nos. 5,476,095; 6,963,771; and 7,361,134; and Pub. Nos. U.S. 2003/0028097; 2006/0094923; and 2009/0236510; and International Pub No. WO 2008/148150 are hereby incorporated by reference for all purposes in their entirety.\nA need exists for a rectal balloon apparatus that relieves rectal gas pressure prior to or during inflation of the balloon. A need also exists for a rectal balloon apparatus that senses the amount of radiation received at or near a treatment site, and that also detects the motion of the rectal balloon apparatus and/or the means used by the apparatus to sense the radiation."}
{"text": "Unproctored (unsupervised) testing of people is suitable in some situations, for example assessing candidates for employment or students for grading purposes. This has occurred because there are some significant cost saving in unproctored testing when compared to proctored (supervised and typically onsite) testing. However the cost saving comes with some disadvantages.\nIn particular, the unproctored testing process is more open to abuse than proctored testing. For example, cheating is made easier by having someone else perform the test on the person's behalf. This is unfair to the remaining candidates or students and at the same time lessens the accuracy and therefore value of the tests for the organisation relying on the testing results."}
{"text": "In recent years, a technology relevant to an antenna for wireless communication, which can be applied to a wearable device (device that a user can wear to use), such as a device of a watch type and a device of a glasses type, has been developed. Technologies relevant to the above antenna are, for example, the technologies described in below Patent Literatures 1 to 3."}
{"text": "Other than referring to a store directory (if one exists) or finding a store sales clerk or a knowledgeable shopper, there is currently nothing available to help shoppers locate items in a store. Thus, many shoppers may search all over a store to locate one or more items, or may wait in a long line at a service counter to ask where particular items are located. This is a particularly difficult problem for handicapped shoppers who use motorized chairs. Such a shopper must press a button, or the like, in order to control movement of the chair. While searching for the location of store items, these shoppers may experience discomfort or fatigue from controlling the motorized chair. Further, shoppers who are hearing impaired may not be able to easily communicate with store employees and may not be able to effectively ask where an item is located. A quick and easy way to locate items in stores is needed to make shopping trips shorter and more productive."}
{"text": "The present invention relates to an image recording apparatus of the the which employs a photo- and pressure-sensitive heat-developable material which is composed of a substrate and a layer formed on the substrate from a material which is photosensitive and which enables a latent image thereon to be developed by heat and the developed image to be fixed by application of pressure. The invention also is concerned with an image recording method which makes use of such a photo- and pressure-sensitive heat-developable material.\nAn image recording material employing microcapsules containing a photosensitive composition has been known. An example of such a recording material is disclosed in Japanese Kokai No. 179,836/1982. This recording material has a substrate which carries capsules made of a synthetic polymeric resin and containing a vinyl compound, a photopolymerization initiator and a coloring precursor.\nIn recording an image, using this recording material, the material is first exposed so that the microcapsules are hardened in conformity with the form of the image. Then, as the recording material is pressed, the microcapsules which have not been hardened are ruptured so as to release the coloring agent precursor, thus forming a color image. This image recording system provides a high quality of the recorded image with a dry-type simple processing, but suffers from a disadvantage in that the photosensitivity is much smaller as compared with the system which makes use of silver halide.\nUnder this circumstance, the present applicant has proposed, in the specification of Japanese Kokai No. 275,742/1986, a novel recording material which has a high photosensitivity and which ensures a high quality of the recorded image with a simple dry-type process. This recording material is a photo- and pressure-sensitive heat-developable material constituted by a carrier and a layer formed on the surface of the carrier, the layer containing at least a photosensitive silver halide, reducing agent, polymerizable compound and a color image forming substance, wherein the polymerizable compound and the color image forming substance are confined in common microcapsules.\nAn image recording method which makes use of this photo- and pressure-sensitive heat-developable material is disclosed in the specification of Japanese Kokai No. 278,849/1986 filed by the same applicant. According to this method, the heat-developable material is first exposed so that a latent image is formed thereon in conformity with the image to be recorded. Then, the material is heated for development so that the polymerizable compound in the area where the latent image exists is polymerized to produce a polymeric compound, thus thermally setting the microcapsules. Then, the material is superposed on an image receiving material having an image receiving layer capable of receiving the color image forming substance, and is pressed to the image receiving material so that at least part of the microcapsules having no latent image is ruptured so as to transfer the color image forming substance to the image receiving material, thereby forming an image on the latter.\nIn order to obtain a clear image by the transfer of the heat-developed image from the photo- and pressure-sensitive heat-developable material to the image receiving material, it is essential that the heat-developable material be pressed uniformly onto the image receiving material under a given condition."}
{"text": "1. Field of the Invention\nThe present invention relates to a capacitor in a semiconductor device and a manufacturing method thereof, and more particularly, to a capacitor including a perovskite series dielectric film and a manufacturing method thereof.\n2. Description of the Related Art\nAs integration of semiconductor devices progresses, the area occupied by capacitors is increasingly reduced. Thus, various approaches have been made to obtain high and stable capacitance in a relatively small area. These approaches include a method of using dielectric films in the metal oxide series such as tantalum oxide (Ta2O5), aluminum oxide (Al2O3), and titanum oxide (TiO2), or dielectric films in the perovskite series such as BST((Ba, Sr) TiO3), SBT((SrBi2Ta2O9), and PZT((Pb,Zr)TiO3). These dielectric films have dielectric constants of tens to hundreds times higher than conventional dielectric films made of conventional silicon nitride or silicon oxide film or composite film thereof have.\nAmong them, using BST as a dielectric film in dynamic random access memory (DRAM) or ferroelectric RAM (FRAM) has been continuously pursued. However, in the case of BST, when it is deposited as a thin film for a capacitor in high integrated semiconductor devices, there is a limitation in reducing the thickness of the film due to poor leakage current characteristics. This limits high integration of semi conductor devices using BST as a dielectric film of a capacitor.\nAccordingly, regardless of the causes of poor leakage characteristics of BST dielectric film, methods of reducing the leakage current of a BST dielectric film capacitor are under study. For example, a method of forming buffer layers has been proposed by U.S. Pat. No. 5,471,364. In this case, the buffer layers are made of low dielectric constant materials formed between a BST dielectric film and an upper electrode or a lower electrode but that exhibit good leakage current characteristics. In addition, the buffer layers are thinner than the BST dielectric film. Furthermore, another method is proposed by U.S. Pat. No. 5,731,220 to reduce leakage currents. In this method, a small amount of erbium (Er) is added to a BST precursor solution with the molar ratio of the erbium to titanium in the BST precursor solution between 0.01-0.05 to deposit a BST film.\nA method of interposing buffer layers, however, has a problem in that forming a dielectric film of a relatively low dielectric constant results in a reduction in the overall capacitance. Furthermore, a method of adding a small amount of Er has many problems in mass production because the process is complicated. Also, changes in capacitance characteristics resulting from changes in process variables are difficult to predict. Additionally, in U.S. Pat. No. 5,731,220, only the importance of reducing the leakage currents by addition of ER has been described, but the implementation of the invention has not been successful because the mechanism for such reduction is not clearly identified.\nAccordingly, there is a need to improve the method of depositing a perovskite series dielectric film, e.g., a BST film, to avoid the problems described above, for example, reduction in overall capacitance, complicated process steps, and unpredictable changes in capacitor dielectric characteristics caused by changes in process conditions, particularly, during actual mass production.\nTo solve the above problems, the present invention provides a capacitor having a perovskite series dielectric film, in which the leakage current characteristics are improved.\nIn addition, the present invention provides a simplified method of manufacturing a perovskite series dielectric capacitor which improves the leakage current characteristics of the dielectric film.\nAccordingly, the present invention reduces leakage currents by containing CuXO or copper in a perovskite series dielectric film. That is, a capacitor according to the present invention includes a perovskite series dielectric film disposed between upper and lower electrodes made of conductive material, wherein the perovskite series dielectric film contains CuXO or copper. The CuXO or copper is preferably contained in the grain boundary of the perovskite series dielectric film.\nMore preferably, the perovskite series dielectric film is a barium strontium titanate (BST) film or a PZT film.\nAlso, the present invention provides a method of manufacturing a perovskite series dielectric film capacitor, including penetrating the perovskite series dielectric film with CuxO or copper. According to a method of manufacturing a general perovskite series dielectric film capacitor, after depositing a conductive material on a substrate to form a lower electrode and depositing a perovskite series dielectric film on the lower electrode, a conductive material is deposited on the perovskite series dielectric film to form an upper electrode. The method of manufacturing the perovskite series dielectric film capacitor according to the invention has steps similar to the method of manufacturing the general perovskite series dielectric film capacitor, further including unique steps of depositing CuxO (x=1 or 2) on the perovskite series dielectric film to form a CuxO film, and performing a heat treatment on the substrate to permit CuxO or copper of the CuxO film to penetrate the perovskite series dielectric film.\nIn this case, the CuXO film may be formed sufficiently thin that substantially all CuXO or copper penetrates the perovskite series dielectric film to leave substantially no residue on the perovskite series dielectric film after the heat treatment. This is for preventing a reduction in the overall capacitance caused by CuxO film having a relatively low dielectric constant compared to the perovskite series film. Furthermore, the heat treatment is performed before, after, or both before and after forming the upper electrode.\nIt is believed that the prevention of leakage current flowing through a grain boundary of a columnar crystal structure of the perovskite series dielectric film is made possible by permitting CuXO or copper to penetrate the grain boundary of the perovskite series dielectric film."}
{"text": "Conditional Random Fields (CRFs) are probabilistic models for computing the probability p(Y|X) of a possible output label sequence Y=(y1, y2, . . . , yn) given the input feature sequence X=(x1, x2, . . . , xn). CRFs are widely used in Natural Language Processing (NLP). For example, CRFs may be used in the frontend of TTS (Text To Speech) to implement word segmentation, POS (Part Of Speech) tagging, prosodic break prediction, and the like. CRFs use an N-gram feature (generally bigram) to model the relationship between neighboring output types. The computation complexity grows exponentially as N increases."}
{"text": "As a method for forming a fine wiring by electroless plating, for example, Patent Literature 1 discloses a method where a resin base surface with excimer light in a desired pattern, the resin base is then brought into contact with a silane coupling agent having an amino group or an amino-modified reactive silicon oil, a receptive layer for trapping/immobilizing a catalyst for initiating an electroless plating reaction is formed on the part irradiated with the excimer light, and a nickel/phosphorus plating layer is formed on the receptive layer by electroless plating."}
{"text": "Much of the earth's remaining natural gas and oil is located in rock shale and rock formations at depths varying from 500-20,000 feet below the surface. Many of the naturally occurring formations suffer from low porosity and permeability, thereby restricting natural gas and/or oil to flow into a well bore where the gas and/or oil can be recovered.\nHydraulic fracturing is a process that creates fractures in low porosity, low permeability rocks of geological formations. A hydraulic fracture may be formed by pumping a fracturing fluid mixture into a well bore at a rate sufficient to increase pressure down-hole to a value in excess of the fracture gradient of the formation rock. The pressure causes the formation to crack or fracture, which allows the fracturing mixture to enter and extend the crack further into the formation.\nTo keep a fracture open after the pumping process stops, the fracturing fluid mixture contains a solid, called a proppant, that remains in the new fracture and keeps the fracture open. Various types of proppant may be used depending on the permeability or grain strength needed. A completed fracture provides a conductive path connecting a large area of the reservoir to the well, thereby increasing the area from which oil, natural gas and liquids can be recovered from the targeted formation.\nPropped hydraulic fracture stimulation is widely used to improve well productivity from tight, otherwise noncommercial reservoirs. Other areas of application include sand control in weakly consolidated reservoirs, gas condensate fields, and high permeability reservoirs that show significant permeability anisotropy.\nAn important consideration in the practice of hydraulic fracturing is “proppant flow-back”, i.e., the production of proppant back to the surface. Proppant flow-back can lead to damaged equipment and downtime, which necessitates costly and manpower intensive surface handling procedures. Proppant flow-back also presents a problem of proppant disposal. In some cases, a well can be prematurely abandoned if costs to return the well to production are excessive. Proppant flow-back, although not desirable, can be tolerated in certain operational environments. However, proppant flow-back from a fracture during the production phase is problematic. In practice, 20% to 50% of the original proppant pumped into a formation may be produced back.\nAn integral part of a hydraulic fracturing process is the introduction of proppant particles of various sizes into a fracture to hold the fractures open. Examples of proppants include structural materials such as sand, man-made ceramics, walnut shells and even polymer beads. Additionally, resin coated proppants have been used in an effort to reduce proppant flow-back. The resin coated proppant approach of gluing the sand together at the interface was an excellent starting point to reduce flow-back but mechanical cycling tends to crack the bonds that interconnect the proppant grains, thereby creating failure.\nStudies have been undertaken to determine the mechanisms of flow-back. Laboratory work and theoretical simulations of several thousand individual proppant grains have shown that flow-back of plain proppant is critically dependent upon the ratio of mean grain diameter to fracture width. This dependence can be visualized as larger “columns” of proppants that have a greater tendency to buckle under loading. As a result of closure stress and inter-particle friction, a geometrically irregular arch of proppant grains is formed behind the fracture mouth. The buckled columns of proppant in front of the arch carry virtually no load. Therefore, small levels of proppant are initially transported to the well-bore even at low fluid flow rates. The loaded arch is stable until a sufficient fluid drag force is reached that is able to collapse the arch. A collapsing arch results in further proppant flow-back. An arch is a curved structure spanning an opening, serving to support a load by resolving the normal stress into lateral stress. For an arch to form, closure stress is required to generate a reaction force with the face of the fracture. The resolution of the forces in an arch necessarily sets up a pattern of shearing stresses within the material. At high closure stress these shear forces can become excessive and lead to pack failure.\nModes of proppant pack failure leading to proppant flow back are thought to be caused by arch destabilization due to excessive hydraulic and/or gravitational forces. Other modes of proppant pack failure are believed to include a fracture closure stress that is too low, proppant pack shear failure caused by too high fracture closure stress, as well as possible proppant crushing at high closure stress.\nCritical parameters determining flow-back of plain proppants include fracture closure stress, both stabilizing and destabilizing. Other critical parameters include hydrodynamic force imposed from fluid production, which is destabilizing as the fluid flow tends to buckle proppant columns at the fracture mouth. Other critical parameters include the fracture width, which affects arch geometry and the transfer of friction and stress forces acting on proppants at a free face. Thus, for the same gradient, larger proppants will experience a greater destabilizing force. Similarly, at the fracture face the larger the proppant the fewer proppants per unit area are available to resist the applied closure stress and the larger the normal stress at the proppant contact.\nA typical fracturing fluid mixture is made up of carrier fluid, additives and proppants. Typically, the carrier fluid and proppant comprise 99% of the mixture with less than 1% additives.\nThere are five types of carrier fluids normally used in the process. The five types are, 1) Water: which is typically un-gelled freshwater or formations of brine; 2) Cross-linked water-based fluids, which are gelled with polymers that employ a cross-linking agent, such as a metallic ion to bond the polymer molecules together for increasing fluid viscosity; 3) Oil-based fluids, which include gelled oils, diesel or lease crude; 4) Oil-in-water emulsions, which are external phase gelled water, internal phase gelled diesel, lease crude or condensates; and 5) Foam, which is a mixture of gas, e.g., nitrogen or carbon dioxide, gelled liquid, such as water or oils, and foaming agents, where mixtures are typically 60% to 80% gas.\nExamples of typical additives used in the hydraulic fracturing industry include: acids; alcohols; bases; biocides; buffers; breakers; clay stabilizers and kcl substitutes; cross-linkers such as aluminum, boron, titanium, or zirconium; cross-link accelerators and delay agents; demulsifiers; foamers and defoamers, friction reducers; iron control agents; cellulose and guar polymers including standard, hpg, cmg, cmhpg; polymer slurries, including diesel and “green”; oxygen scavengers; salts; surfactants; and flurosurfactants. This list of additives is only a representative portion of all additives used and is not meant to be a complete list.\nEffectively dispersing a well-formulated fracturing fluid with a complex set of additives into a homogeneous emulsion is extremely difficult to accomplish due, primarily, because of chemical affinity mismatches that result from the broad diversity of the additives used. Additionally, because of the toxic nature of some of the chemicals used in hydraulic fracturing, industry is facing stiff and growing environmental regulation. Regulation is anticipated that will establish limits for various materials contained in hydraulic fracturing additives and carrier fluids.\nDue to environmental issues associated with hydraulic fracturing, the oil and gas industry has shifted towards using water with mineral additives as the fracturing fluid of choice. The most commonly used additive in water hydraulic fracturing is a friction reducer.\nCarrier fluids, additives and proppant must be mixed for use in hydraulic fracturing.\nThere are three types of commonly used mixing principles:\n1. Static mixing, wherein liquids flow around fixed objects, either via force produced flow by pressure through mechanical means or gravity induced flow;\n2. Dynamic mixing, wherein liquid induced mixing results from mechanical agitation via impellers, wiping blade and high shear turbines as well as single or double screw extruder designs or screw agitation designs; and\n3. Kinetic mixing, wherein liquid is mixed by velocity impacts on a surface or wherein two or more liquids form a velocity impact by impinging on one another.\nAll three of the above mixing methods have one thing in common that hinders the optimizing of mixing regardless of the fluid being combined and regardless of whether the materials being mixed are polar, nonpolar, organic or inorganic etc. or if it is a filled material with compressible or non-compressible fillers.\nThe commonality that hinders optimizing of mixing is that all incompressible fluids have a wall effect or a boundary layer effect where the fluid velocity is greatly reduced at the wall or mechanical interface. Static mixing systems use this boundary layer to fold or blend the liquid using this resistive force to promote agitation.\nDynamic mixing, regardless of the geometry of mixing blades or turbine, results in dead zones and incomplete mixing because of the boundary layer. Dynamic mixing uses high shear and a screw blade designed to use the boundary layer to promote friction and compression by centrifugal forces to accomplish agitation while maintaining an incomplete mixed boundary layer on mechanical surfaces.\nKinetic mixing suffers from boundary layer effects on velocity profiles both on the incoming streams and at the injector tip. However, kinetic mixing suffers minimal boundary layer effects except for transport fluid phenomena.\nA further explanation of the boundary layer of the flowing fluids follows. Aerodynamic forces depend in a complex way on the viscosity of a fluid. As a fluid moves past an object, the molecules in close proximity to the surface tend to stick to the surface. The molecules just above the surface are slowed down by their collisions with the molecules that are sticking to the surface. These molecules, in turn, slow down the flow just above them. The farther away from the surface, the fewer the collisions that are affected by the object surface. Therefore, a thin layer of fluid is created near the surface in which the velocity of the fluid changes from zero at the surface to the free stream value away from the surface. Engineers call this layer the “boundary layer” because the layer occurs on the boundary of the fluid.\nAs discussed above, as an object moves through a fluid, or as a fluid moves past an object, molecules of the fluid near the object are disturbed and as the molecules move around the object. Aerodynamic forces are generated between the fluid and the object. The magnitude of the aerodynamic forces depends on the shape of the object, the speed of the object, the mass of the fluid going by the object and on two other important properties of the fluid, i.e., the viscosity, or stickiness, and the compressibility, or springiness, of the fluid. To properly model these effects, aerospace engineers use similarity parameters, which are ratios of these effects to other forces present in the problem. If two experiments have the same values for the similarity parameters, then the relative importance of the forces are being correctly modeled.\nReferring now to FIG. 1, a two dimensional representation of the streamwise velocity variation from free stream to the surface is shown. In reality, the effects are three dimensional. From the conservation of mass in three dimensions, a change in velocity in the streamwise direction causes a change in velocity in the other directions as well. A small component of velocity perpendicular to the surface displaces or moves the flow above it. The thickness of the boundary layer can be defined as the amount of this displacement. The displacement thickness depends on the Reynolds number, which is the ratio of inertial (resistant to change or motion) forces to viscous (heavy and gluey) forces and is given by the equation: Reynolds number (Re) equals velocity (V) times density (r) times a characteristic length (l) divided by the viscosity coefficient (mμ), i.e., Re=V*r*l/mμ.\nStill referring to FIG. 1, boundary layers may be either laminar (layered), or turbulent (disordered) depending on the value of the Reynolds number. For lower Reynolds numbers, the boundary layer is laminar and the streamwise velocity changes uniformly as a function of distance away from the wall, as may be seen on the left side of FIG. 1. For higher Reynolds numbers, the boundary layer is turbulent and the streamwise velocity is characterized by unsteady, or changing with time, swirling flows inside the boundary layer. The external flow reacts to the edge of the boundary layer just as it would to the physical surface of an object. Therefore, the boundary layer gives any object an effective shape, which is usually slightly different from the physical shape. The boundary layer may lift off or “separate” from the body and create an effective shape that is substantially different from the physical shape. Separation occurs because the flow in the boundary has very low energy relative to the free stream and is, therefore, more easily driven by changes in pressure. Flow separation is the reason for airplane wing stall at high angle of attack.\nBoundary-layer Flow\nThe portion of a fluid flow that occurs near a solid surface is where shear stresses are significant and is where inviscid-flow assumptions may not be used. Solid surfaces interact with a viscous fluid flow because of the no-slip condition discussed above, i.e., because of the physical requirement that fluid and solid have equal velocities at their interface. Therefore, fluid flow is retarded by a fixed solid surface, which results in the formation of a finite, slow-moving boundary layer. For the boundary layer to be thin, the Reynolds number of the body must be large, i.e., 103 or greater. Under these conditions the flow outside the boundary layer is essentially inviscid and plays the role of being a driving mechanism for the layer.\nReferring now to FIG. 2, a typical low-speed or laminar boundary layer is shown. Such a display of the streamwise flow vector variation near a wall is called a velocity profile. The no-slip condition requires that u(x,0)=0, as shown, where u is the velocity of flow in the boundary layer. Velocity rises monotonically with distance y from the wall, finally merging smoothly with the outer (inviscid) stream velocity U(x). Assuming a Newtonian fluid, at any point in the boundary layer the fluid shear stress τ is proportional to the local velocity gradient. The value of the shear stress at the wall is most important, since the shear stress value relates not only to the drag of the body but often also to its heat transfer. At the edge of the boundary layer τ approaches zero asymptotically. There is no exact spot where τ=0, therefore the thickness δ of a boundary layer is usually defined arbitrarily as the point where u=0.99 U.\nFriction Reducers\nAs stated above, additives in hydraulic fracturing fluid are typically deployed for use as friction reducers. As shown in FIG. 3, experiments were conducted wherein various friction reducers were added at a concentration of 0.25 gpt in 2% (wt) Kcl tap water flowing through a ½″ OD/0.402″ ID pipe. FIG. 4 shows experimental results comparing of 1.0 gpt of cationic friction reducers in several different fluids through a ½″ OD/0.402″ ID pipe. The graphs are discussed in SPE 119900, “Critical Evaluations of Additives Used in Shale Slickwater Fracs”; (P. Kaufman and G. S. Penny, CESI Chemical a Flotek Co. and J. Paktinat, Universal Well Services Inc.)\nEach of FIGS. 3 and 4 show that there are two clear regions where the dispersion of a typical material is greatly influenced, i.e., at 0 to 10 seconds and again at 20 to 30 seconds. During the experiments, the fluid was recycled back into the loop through the circulation pump every 10 seconds. The circulation pump acted like a mixing system, thereby improving the material performance of the additive. The experimental results, represented by the curves of FIGS. 3 and 4, illustrate that dispersion of additives is crucial to the performance of the material. The periods of time from 0 to 10 second and from 20 to 30 seconds correlate to the times when typical friction reduction material passes through the mixing system in the cyclic loop process. The shape of the curves of FIGS. 3 and 4 indicate that, at about 10 seconds, the typical material has not been adequately dispersed and, therefore, is not distributed into the boundary layer of the flowing material. Therefore, the important disbursement of additives is desirable."}
{"text": "The invention relates to a method of producing an inspection chip used for a specimen detection device that detects a specimen using localized plasmons and a specimen detecting method of detecting a specimen using localized plasmon.\nAmong known detection devices for detecting a minute specimen is one using an enhanced electric field created by localized plasmons or an enhanced electric field further intensified by localized plasmon resonance, which is a phenomenon where resonance is produced by coupling between localized plasmons and light.\nIn a detection device of this type, a specimen is placed on a substrate on a surface of which a number of metallic portions capable of exciting localized plasmons are disposed, with Raman scattered light of the specimen intensified by an enhanced electric field that is generated by localized plasmons. Since Raman scattered light displays a spectrum proper to each specimen under examination, the specimen can be identified or existence of a specific specimen can be determined by detecting Raman scattered light of the specimen.\nSubstrates that generate an enhanced electric field intensified by localized plasmon resonance include, for example, ones having configurations described in JP 2005-144569 A and JP 2007-240361 A.\nJP 2005-144569 A describes a substrate on which particles each coated with a thin metallic film are disposed with a high density. To be more specific, JP 2005-144569 A describes a substrate having a two-dimensional array structure formed by reducing dielectric or semiconductor particles densely arranged on the substrate through anisotropic dry etching process and attaching a metal or a semiconductor onto the particles in a hemispherical form by vapor deposition or sputtering to keep the metallic parts of adjacent particles spaced at a given controlled distance.\nJP 2007-240361 A describes a substrate having a number of minute projections formed thereon. The projections have a metallic layer disposed on the surface thereof; the metallic layer has a dielectric layer formed on the surface thereof, and the dielectric layer has secured to the surface thereof specific binding members capable of specifically binding to a substance to be measured that is contained in a test sample to form a specific binding substance.\nJP 2005-195440 A describes a specimen detecting method using a microstructure comprising a substrate and first metallic particles. The substrate has micropores formed and distributed on one surface thereof. The first metallic particles each have dimensions permitting excitation of localized plasmon resonance and are disposed in the micropores of the substrate so that the head portions of the first metallic particles project above the surface of the substrate. The microstructure is placed in a solution with a substance capable of specifically binding to the specimen secured to the head portions of the first metallic particles. The solution contains dispersed therein second metallic particles having dimensions permitting excitation of localized plasmon resonance, with a substance capable of specifically binding to the specimen secured to the second metallic particles. The surface of the microstructure from which the head portions of the first metallic particles project is irradiated by light to measure the intensity of light component reflected or transmitted at said surface, thereby detecting the specimen contained in the solution based upon the measured intensity of the light.\nAccording to the method described in JP 2005-195440 A, the specimen binds to the second metallic particles to which a substance capable of specifically binding to the specimen is secured whereas the specimen also binds to the first metallic particles to the head portions of which a substance capable of specifically binding to the specimen is secured, so that the change in plasmon resonance wavelength that varies greatly when the first metallic particles and the second metallic particles approach is detected to achieve detection of the specimen.\nIn the detection devices using the substrates described in JP 2005-144569 A and JP 2007-240361 A, the metallic particles or metallic projections having dimensions permitting excitation of localized plasmon resonance are allowed to adsorb specimen to detect, for example, the variation in an enhanced electric field, fluorescent light from the specimen excited by the enhanced electric field, and Raman scattered light thereby to achieve detection of the specimen.\nAccording to such a method, however, a great enhancement cannot be achieved because only an enhanced electric field created around metallic particles that lie below the specimen can be used.\nAlthough, according to the method described in JP 2007-240361 A, the efficiency with which the specimen binds to the metallic film may be increased by securing a substance capable of specifically adsorbing the specimen to the surface of the metallic film, a great enhancement cannot be achieved because, also in this case, only an enhanced electric field created around metallic particles that lie below the specimen can be used.\nFurthermore, the need to select for each specimen a specific adsorbing substance for attachment to the metallic particles adds to the cumbersome procedure and limits the specimen to which the detection device can be used.\nThe method proposed in JP 2005-195440 A, which detects the variation of a plasmon resonance wavelength, requires the specific adsorption to be performed twice. Accordingly, two or more molecules of a specific adsorbing substance (referred to as “antibody” below) needs to be bound to the specimen. Thus, where the specimen is a substance having a small molecular mass, two molecules of an antibody cannot be adsorbed onto the specimen, making the detection of the specimen impossible.\nFurthermore, because the method described in JP 2005-195440 A achieves detection of the variation in plasmon resonance wavelength caused as the first metallic particles and the second metallic particles come close to each other through the intermediary of the specimen, it is essential that the first metallic particles and the second metallic particles approach each other only through the intermediary of the specimen."}
{"text": "The present invention relates generally to signal spectrum analyzers and more particularly, to such analyzers wherein Chirp-Z forward transforms are utilized.\nConventional Chirp-Z signal analyzers utilize forward transforms to compress the individual continuous wave (hereinafter CW) components of the signal into segregated pulses or time intervals from which such components are then selectively reconstructed or expanded with inverse transforms. Considerable gain is inherent with Chirp-Z forward transforms, however, high attenuation is inherent with Chirp-Z inverse transforms. Consequently, Chirp-Z inverse transforms encumber the overall sensitivity of signal analyzers."}
{"text": "Field of the Invention\nThe present invention relates to a clock such as an alarm clock, a wall clock, or a watch to be mounted on a knapsack and similar articles. More specifically, the invention relates to interchangeable decorative members which may be mounted on the alarm clock, wall clock or watch to be inserted in the knapsack or similar articles."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for the determination of the amount of irreversibly glycated hemoglobin, or HbA1c, present in a sample of blood, relative to the amount of total hemoglobin. In particular, the invention incorporates in the method an electrochemical, enzyme-catalyzed reaction or reactions.\n2. Background Information\nHbA1c is a glycated hemoglobin formed by a binding reaction between an amine group of hemoglobin and the glucose aldehyde group, for example between the amine group of the N-terminal valine of the β-chain of hemoglobin and the glucose aldehyde group. The binding reaction first forms a Schiff's base and then a stable ketoamine by Amadori rearrangement. The percentage of HbA1c (i.e. the amount of glycated hemoglobin relative to total hemoglobin in the blood) has come to be taken as a measure of the level of blood glucose control a diabetic patient has maintained for a period of two or three months prior to the measurement. As such, percentage HbA1c has become an important measurement by which health care providers can assist diabetic patients in their care.\nThere are many known assays that can be used to determine HbA1c percentage. In recent years research efforts have focused on creating assays that are both highly accurate and fast. However, known HbA1c assays typically require a substantial number of time-consuming steps wherein the blood components must be separated and treated.\nIn the health care context, a diabetic patient is typically guided by a physician to obtain an HbA1c measurement when the physician realizes that there is a need for such information during an office visit. The patient then provides a blood sample to a laboratory and results are returned to the physician hours or days later. Typically, the lab will use a table top analyzer of the type presently available commercially. This time lag between the patient's visit and the result of the test requires that the physician review the result long after the patient has left the office. If the physician believes that further consultation with the patient is required in light of the test result, the patient must be contacted again.\nCurrently, there is a device sold under the name “A1c NOW” by Metrika, Inc. of Sunnyvale, Calif. This handheld and disposable device (based on technology described in U.S. Pat. No. 5,837,546 entitled “Electronic Assay Device and Method,” incorporated herein by reference) is said to provide an HbA1c test result in eight minutes using a relatively small sample of blood. The A1c NOW device is an example of the market demand for a fast method of providing an HbA1c result for either home or doctor's office use. However, the A1c NOW device is not as accurate as some laboratory assays. Thus, research has continued to focus on finding a highly accurate HbA1c assay that is also fast enough and simple enough to permit a diabetic and his or her doctor to take a blood sample during an office visit and have a trustworthy HbA1c measurement available for discussion in the same visit."}
{"text": "Light emitting diodes (“LEDs”) have been used in various applications including illuminating watches, transmitting information from remote controls, and forming images on jumbo television screens. More recently, LEDs have been used in portable lighting devices, such as flashlights, because, among other things, LEDs can last longer and can be more durable than incandescent lamps commonly used in conventional flashlights.\nNotwithstanding the desirable characteristics LEDs may have over incandescent lamps, improvements can be made over existing lighting devices that use an LED as its primary source of light. For example, current LED flashlights typically fail to produce a quality light beam that projects for any appreciable distance. A reason for this is that available LED lamps substantially radiate light in a pattern over an angle less than 180° relative to the position of the LED. FIG. 1 illustrates the brightness or radiation pattern of a typical LED lamp. A typical LED lamp includes an LED and a lens arranged over the LED. The light rays that emanate from the typical LED lamp is generally conical —the brightness (indicated in percentages) is generally concentrated about the center axis 11 and reduces non-linearly as the spherical angle θ increases. Existing LED flashlights have not provided a reflector/lamp combination that effectively captures the brighter light rays that are concentrated about the center axis. Accordingly, although available LED lighting devices may be suitable to illuminate the immediate surrounding area, the distance that the light beam is able to project has been limited.\nTo try to overcome this deficiency, some devices have used multiple LEDs or a combination of LED lamps and incandescent lamps. However, such devices involve greater complexity, consume more energy, and cost more to manufacture. Accordingly, the present invention provides an energy efficient LED lighting device that has improved optical performance and that projects a quality light beam.\nAlso, as improvements are made to light sources, such as LEDs, another problem challenging the operation of portable lighting devices is effectively dissipating the increased heat that is generated by the light source. Accordingly, the present invention provides a combination that effectively dissipates heat from the light source of a portable lighting device. The present invention also provides a means to reduce the amount of heat generated by the light source and to use less energy to illuminate the light source."}
{"text": "Field of the Invention\nThe present invention relates to methods, systems, apparatuses, computer readable media, and other means for reliably using thermophotovoltaic cells to power electronic devices.\nDescription of Related Art\nElectronic devices are known which are powered by rechargeable batteries such as nickel-cadmium, lithium-ion, nickel-metal hydride, and rechargeable alkaline batteries. The batteries of such devices are often recharged using standard recharging means and methods, such as by plugging a charging cable into the device and into a conventional alternating current (AC) outlet. The battery of a device can be conventionally charged by plugging the device into another electronic device using a universal serial bus (USB) connection (such as a USB cable or docking station). However, it can be inconvenient or impossible for the user to charge the batteries when, for example, the user is traveling. In addition, a typical recharge will power the battery for only a limited amount of run time. After the battery has been drained, the user must recharge or replace the battery to continue using the electronic device.\nSome electronic devices (e.g., calculators) can further use solar cells to power the device. However, such methods and devices are limited by small size of many electronic devices, which constrains the effective surface area for absorbing sunlight, as well as the inherent characteristics of solar cells, which require the device to track the sun to maintain consistent exposure to sunlight for stable power output. Consequently, while the voltage generated by the solar cells can be enough to power less energy demanding devices, the voltage can be insufficient for electronic devices that demand more power. Furthermore, as electronic devices include an increasing number of advanced features, larger sources of stable and consistent power are required. Thus, there are competing interests between the portability of these devices and the amount of available power.\nThe effectiveness of using known photovoltaic cells to charge electronic devices is additionally limited by composition and other properties of the cells which frustrate charge separation and transport and permit significant electron-hole pair recombination to occur. Unless rapidly separated and transported from the cell, photogenerated charges recombine and contribute to thermalization and, in low band gap semiconductors, photocorrosion. Increasing the thermal energy of such cells further reduces overall photoconversion efficiency (PCE) by lowering charge mobility and promoting additional recombination.\nThermalization inevitably occurs in such applications notwithstanding the relative speed and degree of separation and transport. In response to AM 1.5 G solar radiation, crystalline silicon cells with a 1.1 eV absorption edge are able to absorb photons representing approximately 77% of the solar spectrum (with wavelengths ranging from about 350 nm to about 1.1 μm); the remainder (with wavelengths ranging from about 1.1 μm to 1 mm) heats the cell and its surrounding structures and environment. Further, almost all of the photons in the flux absorbed by the cell have energies greater than the bandgap energy of silicon, and produce both an electron-hole pair and thermalization losses upon absorption. In total, more than about two thirds of the solar flux striking a conventional solar cell is converted to thermal energy.\nProblematically, the performance of electronic devices also degrades in response to elevated temperature. The integration of conventional solar cells into electronic devices is thus further complicated by the implicit expectation that such devices be exposed to solar radiation. However, electronic devices generate heat since such devices are comprised in part of components which only use a portion of the charge they receive. For example, the amount of power discharged from a device's battery that is not used for device function instead increases the thermal energy of the device at rates dependent on the relative efficiency of the device's components. Thus, device manufacturers go to significant lengths to manage heat, in many cases by designing their devices to dissipate heat as a form of waste since it is generally present at low grade temperatures insufficient for known heat recovery technologies to operate cost-effectively, particularly at the scale of a electronic device.\nThe thermal energy produced by or imparted to known electronic devices, be it from the absorption and inefficient conversion of sunlight or the inefficient use of electrical energy, will eventually dissipate as the system and its surroundings achieve thermal equilibrium. The rate and degree to which that dissipation occurs is in part a function of the ambient temperature in the local environment, which is itself directly and indirectly heated by the vast majority of the incoming solar flux that never comes in contact with the surface of the cell. Most of that energy is absorbed and stored by the environment (e.g., air, water, land, buildings, electronic devices, and so on) in the form of thermal energy which is then released at low grade temperatures by convection, conduction and, importantly, emission of infrared radiation in every direction, day and night, 24 hours a day, and 365 days per year. Humanity then adds to those emissions by engaging in activities that create waste heat, for example, by burning coal and other fossil fuels to provide power to electronic devices. Historical attempts to access and convert such low grade sources of thermal energy into electricity have been limited and unsuccessful.\nThe increased use of electronic devices in recent years corresponds to increased fossil fuel consumption. For example, the two billion smartphones expected to be in use worldwide by 2014 will consume the equivalent of more than 80 million barrels of fossil fuel per year. Smartphone use alone is growing by more than about 40% per year, with each new generation evolving and selecting for increased performance. It is desirable to provide a method or device for meeting the power needs of such devices that does not rely on the combustion of fossil fuels.\nIn view of the conventional attempts, needs remain for improved methods, devices and systems for powering and managing thermal energy in electronic devices."}
{"text": "In the past, if solar energy concentrating systems were used on top of buildings or roofed structures to gather radiant solar energy, then one could select from two options. The first option was simply to take a solar energy concentrating system suitable for use on the ground and mount it up on top of a roofed structure. While such an approach did gather radiant solar energy, it also required that the roof support a substantially greater weight. This add-on approach resulted in substantially higher capital costs. Moreover, once the system was atop the roofed structure, inevitably the underlying roof would need repair. The presence of the system atop the roof in a non-integrated fashion increased the ease and cost of making roof repairs.\nThe second option was to make a substantial portion of the roof into a solar energy concentrating reflector. An example of this unitary reflector approach can be found in U.S. Pat. No. 3,994,435 to Barr. While Barr no longer had the disadvantage of added supporting structure, Barr had to make compromises in gathering the radiant solar energy. The semi-cylindrical reflector and fixed collector did not collect as much solar energy as ground-based units with better geometries. Also, Barr required that the underlying building have a shape similar to the reflector, and thus, the ends of the Barr building had to swoop arcuately upwards, mimicking the reflector arc. Such a requirement had obvious disadvantages in being used on the numerous flat roofed buildings and roofed structures which are used for commercial or industrial purposes."}
{"text": "The invention relates to a process for the production of xylenes by adsorption-reaction in a simulated moving bed and particularly a simultaneous process for dismutation and separation of a toluene feedstock into benzene and xylenes.\nParaxylene is a raw material that makes it possible to synthesize terephthalic acid and toluene terephthalate that are then transformed into resins and fibers, such as, for example, Tergal. The most effective process for the production and purification of paraxylene in terms of productivity is the adsorption in a simulated moving bed or countercurrent that may or may not be combined with a final purification by crystallization and necessarily coupled to a catalytic isomerization on a zeolite, for example an EUO-structural-type zeolite that is mentioned, for example, in the French patent application of the applicant 99/07 968. The paraxylene is produced mixed with the other isomers into C8-aromatic compounds primarily by the processes of isomerization, reforming and dismutation of toluene.\nThe dismutation of toluene consists in using two toluene molecules to produce a benzene molecule and a xylene molecule. This dismutation is referred to as paraselective when the paraxylene constitutes the majority of the xylenes that are produced.\nThe non-selective reaction is catalyzed by MFI-type zeolites (or pentasil) or mordenite or else metals of group VIII; examples will be found in Patents FR 2761904, U.S. Pat. Nos. 4,007,231, 4,011,276, 4,029,716 and 4,052,476. In the case of the selective dismutation, the catalysts that are used are pentasils of type ZSM5, ZSM 11, ZSM22 or ZSM23 or ZSM 35 on which a deposit of outside material such as carbon (coke) or silicon was placed, as described in U.S. Pat. Nos. 4,260,843, 4,274,982, 4,380,685, 4,908,342, 5,173,461 and EP26962. The Mobil Oil Company uses two commercial processes for dismutation of toluene: the MTDP (non-selective) (see Oil and Gas Journal Vol. 69 No. 48, 1971 by Grandio et al.), and the MSTDP (selective) (see Oil and Gas Journal Vol. 90 No. 41, 1992 by Gora et al. or Europ. Chem. News Vol. 54, p. 1418, 1990). For these two processes, the major drawback is that the conversion of toluene is limited to 30% by the thermodynamic equilibrium. At the outlet of the reactor, the typical composition of the mixture is therefore benzene 15%, toluene 70%, and xylenes 15%. A series of two columns to be distilled is therefore necessary to separate these three products and to recycle the largest portion of the toluene to the inlet of the reactor. In the case of the paraselective process, among the isomers of the xylene, it is possible to obtain up to 90% of paraxylene, but if a compromise between the conversion of toluene and the paraselectivity is sought, the proportion of paraxylene is smaller (from 75 to 85%).\nSince the 1970\"\"s, it has been known to use adsorption and reaction simultaneously (IandEC Fundam, 17 (1978) 1 D Sweich, J Villermaux). In the 1980\"\"s, it was proposed by Barker and Ganestsos to use the adsorption reaction in a simulated moving bed and in liquid phase, for example, for the inversion of saccharose and its total separation into glucose and fructose. Since the beginning of the 1990\"\"s, the idea of using adsorption and vapor-phase reaction simultaneously by using hydrogen as a vector fluid and as a reagent was first conceived by A. K. Ray and R. W. Carr (Chem Eng. Sci 50 (1995) n14, 2195-2202). It applied to the reaction for hydrogenation of mesitylene. Other authors adopted the idea of the separation reaction in a simulated moving bed and applied it to refining or to petrochemistry. It is thus possible to cite U.S. Pat. No. 5,530,173 that offers the additional original feature that the feedstock (mixture of aliphatic hydrocarbons with 6 to 8 carbon atoms) and the desorbent (normal pentane) are converted and separated simultaneously into isoparaffins with a high octane rating. Even more recently, U.S. Pat. No. 5,877,373 proposed the transalkylation of aromatic hydrocarbons into C9 and C10 by a mixture of toluene and benzene in a vapor phase and in the presence of hydrogen in such a way as to produce xylenes. In the latter document, it is not the implementation in a simulated moving bed but the circulation in opposite directions in the two reactors.\nThe object of the invention is the simultaneous production of benzene and xylenes by dismutation of toluene substantially without recycling toluene at the inlet of the reactor. To do this, the invention uses separation by adsorption, and the simultaneous reaction used in a simulated moving bed and in vapor phase.\nMore specifically, the invention relates to a simultaneous process for dismutation and separation of a feedstock that essentially consists of toluene, benzene and xylenes in the presence of a desorbent that is high in hydrogen in an adsorber-reactor that comprises at least three zones, characterized in that\na) the vapor-phase or supercritical feedstock is introduced at the inlet of a reaction and adsorption zone (zone III), and at the outlet of said zone, a raffinate that is high in benzene and desorbent is recovered,\nb) the desorbent is introduced at the inlet of a desorption zone (zone I), and an effluent is recovered at the outlet of said zone from which a portion in the form of an extract that is high in xylenes and desorbent is drawn off,\nc) the other portion of the effluent of zone I is introduced at the inlet of a reaction and desorption zone (zone II), and an effluent that is sent to the inlet of reaction and adsorption zone III is recovered at the outlet of said zone.\nThe process is also characterized in that each of the zones comprises at least one bed that contains an adsorbent that is adapted to separate benzene, toluene and xylenes and a catalyst that is adapted to dismutate toluene into benzene and xylenes, whereby the adsorbent and the catalyst are in solid form.\nWhen one of the zones comprises several beds, the latter comprises at least one bed that contains an adsorbent and a catalyst, whereby the remaining beds can contain either only a catalyst or only an adsorbent.\nAccording to a variant of the invention, the catalyst can be a non-selective catalyst for paraxylene, adapted to obtain in the extract a mixture of xylenes that is close to the thermodynamic equilibrium.\nIt may be selected from the group that is formed by MFI zeolites or pentasils, mordenite in acid form or a precious metal of group VIIIB that is deposited on mordenite, or the Y zeolite that is exchanged with nickel, sodium, lanthanum or in acid form, the X zeolite that is exchanged with sodium, lanthanum or in acid form.\nAccording to another variant of the invention, the catalyst is selective for the paraxylene and adapted to obtain a mixture of xylenes containing at least 80% of paraxylene.\nThese catalysts have very few active sites on their surface; they can be selected from the group of MFI zeolites, in particular ZSM5, ZSM11, ZSM22, ZSM23 that are made selective by deposition of coke, or silicon, or magnesium, or germanium or a combination of these elements at the surface of the zeolite.\nWith regard to the adsorbent, it may be selected from the group that is formed by silicalite, MFI zeolites or pentasils, faujasites (X or Y zeolites) that are exchanged by a cation of group IA or a cation of group IIA, or a cation of group IA and a cation of group IIA, the mordenites that are exchanged by a cation of group IA or a cation of group IIA, and the preceding zeolites made selective by deposition at the surface of carbon or magnesium or germanium or silicon or a combination of these elements.\nAccording to a first variant, it is possible to use the same solid both as adsorbent and as catalyst in each bed.\nAccording to a second variant, it is possible to produce an approximately homogenous mixture of adsorbent and catalyst in each of the beds in proportions of 95:5 to 5:95.\nAccording to a third variant, it is possible to use a layer of catalyst and a layer of adsorbent in each bed in proportions of 95:5 to 5:95.\nThe operating conditions in the adsorber-reactor are generally as follows:\nthe temperature is between 220 and 520xc2x0 C. and preferably between 275 and 350xc2x0 C.\nthe H2/feedstock molar ratio is in, for example, proportions of 0.5/1 to 50/1 and preferably between 2/1 and 10/1\nthe total pressure in the adsorber-reactor is between atmospheric pressure and 400 bar and preferably between 1.5 and 40 bar (1 bar=105 Pa)\nthe partial pressure of hydrocarbons at the injection point of the feedstock is generally between 0.1 and 40 bar and preferably between 0.2 and 10 bar\nthe volumetric flow rate is between a pph of 0.025 hxe2x88x921 and 25 hxe2x88x921, and preferably between 0.05 and 5 hxe2x88x921.\nAccording to a first implementation of the process according to the invention that corresponds to a 3-zone adsorber-reactor, the fluid circulation between the last bed of zone III and the first bed of zone 1 can be interrupted with an all-or-nothing valve that is placed downstream from the sampling of the entire raffinate.\nAccording to a second implementation of the process that corresponds to an adsorber-reactor with at least 4 zones, a recycling compressor that is located between the first bed and the last bed is adapted to control the flow of fluid that circulates successively in each of the zones.\nAccording to a third implementation of the process that corresponds to an adsorber-reactor with at least four zones, the circulation of fluid between the last bed of zone IV and the first bed of zone I is interrupted with an all-or-nothing valve, and the fluid that contains the hydrogen that is obtained from zone IV is recovered and sent back to the intake of the hydrogen compressor that delivers this desorbent to the inlet of zone I."}
{"text": "The present invention relates to a method and apparatus for the precision cutting of corrugated paperboard panels into specimens having square, or other rectangular, shapes of predetermined size, e.g., four inches on a side, etc. The precision, cut to size specimens are subsequently used as test specimens for various measurements, e.g., bending tests, column strength tests, edge compression tests, flat crush tests, and adhesion tests.\nApparatus has previously been developed for cutting corrugated paperboard panels into square or rectangular configurations for subsequent use as test specimens in standardized testing of the paperboard material. The conventional cutting apparatus can take various forms. One known cutting apparatus comprises a horizontal table or platform having a flat upper surface adapted to receive the corrugated paperboard panel flatwise thereon, and an elongated knife blade swingably connected to the table for downward motion along one edge of the table. A handle extending from the elongated knife is pulled downwardly so that the knife blade slices through the paperboard panel in a downward arc, to thereby form a straight edge on the panel. The panel is then repositioned on the table, and the cutting process is repeated to form four straight edges on the paperboard panel.\nThe above table is equipped with an upstanding linear abutment, or backstop, oriented at right angles to the cutting plane of the knife blade, whereby the paperboard panel can be positioned with a cut straight edge engaged against the abutment, while the various cuts are being made in the panel. The abutment, or backstop, permits successive cuts to be made at right angles to one another, so that the final test specimen has a rectangular, usually square, or other rectangular, configuration.\nThe cutting table, or platform, may have a slidably adjustable T-square, that includes a horizontal upstanding straight edge extending parallel to the cutting plane of the vertically-oriented knife blade. The horizontal spacing between the upstanding straight edge and the blade cutting plane, can be varied or changed to permit variation in the size of the cut specimen, or test specimen panel. As previously noted, the cut test specimen will usually have a square, or other rectangular, shape. Typically, the cut square or rectangular specimen, will vary in size from an upper limit of about six inches along each edge, to a lower limit of about one inch along each edge. The size of the cut test specimen is determined by tests to be performed, and effectuated by the settings of the slidable T-square, or, more properly, an L-square.\nThe cutting blade can be arranged for vertical motion along an edge of the panel support table, as previously described. However, a vertically moving cutter blade does not achieve an optimum cutting action, since it exerts a downward pressure on the upper surface of the paperboard panel. The cut edge may thus be slightly depressed, or squashed, especially if the cutter blade is slightly worn, in which case, there is a greater cutter blade edge area in contact with the paperboard surface.\nAs an alternative, and as practiced in the present invention, the cutter blade can be made to move horizontally through the paperboard panel parallel to the panel faces, i.e., the cutting blade can be arranged to move horizontally along the table surface so as to produce essentially no deformation of the panel face. In such a case, the lower end of the cutter blade extends into a slot in the table surface, while the cutting edge of the cutter blade extends generally vertically, so that the cutter blade attacks the edge of the paperboard panel along a horizontal action line.\nThe cutter mechanism can be operated manually, or with a power actuated mechanism, e.g., a fluid cylinder. In one known apparatus, an air-operated fluid cylinder is supported for horizontal motion across the table surface. A cutter blade is carried by the movable portion of the cylinder, such that pressurization of the cylinder sends the cutter blade across the table to form a cut straight edge on the corrugated paperboard panel. Reverse motion of the cylinder returns the cutter blade to its original position ready for the next cutting cycle, after repositionment of the paperboard panel.\nOne major and serious disadvantage of all existing prior art cutting apparatus, is the fact that there is no mechanism for precisely determining, and effecting, the precise orientations of the test specimen panel cut edges, relative to the internal flutes in the corrugated paperboard panels, prior to making the cuts. The term, \"flutes,\" is herein used to mean the linear passages, or `tunnels,` formed within the corrugated paperboard panel, by the corrugated internal liner sheet.\nTypically, the corrugated paperboard panels have two flat spaced-apart outer sheets, and an internal corrugated liner sheet. The corrugations, or undulations, defining the flutes, space the two outer sheets apart, and determine the overall thickness of the panel. Typically the paperboard panel will have a thickness of about one-eighth of an inch, such that each undulation has a transverse dimension slightly less than about one-eighth of an inch.\nThe square or rectangular test panels are designed to be subjected to several different tests, including an extremely important edgewise compression test, which should be performed `perfectly`, or precisely, parallel, to the direction of the internal flutes, or, alternately, `perfectly` normal to the flute direction. In order to produce reproducible, or duplicate, test results on similarly constructed paperboard panels, it is necessary that the square or rectangular test specimen panels be uniformly and precisely cut, or formed, so that two of the test specimen panel edges are always precisely parallel to the flutes, and the other two specimen panel edges are always precisely normal to the flutes.\nThe hollow, internal flutes are located within the corrugated paperboard panels, so that the flute direction is not always readily determined by a visual inspection of the uncut corrugated paperboard panel. However, it is possible to form a panel edge running \"approximately\" parallel to the flute direction, by a trial-and-error process, wherein the edge of the test specimen, or sample, is inspected after the first cut, to determine the approximate extent of the deviation from parallelism.\nHowever, such visual trial-and-error inspection, involves an inspection of the corrugated panel sheet edges that are in semi-concealed locations between the outer facing sheets of the corrugated paperboard panel. Therefore, the observer is never quite certain, from a visual inspection, whether or not, the first cut is, in fact, exactly, or precisely, parallel to the internal flute direction. In any event, such visual inspections, as flawed as they are, also require a lengthy amount of time, and a certain degree of skill, on the part of the operator.\nIt is possible to pry apart the edges of the outer facing sheets of the paperboard panel, in order to better see the directions of the flutes. However, such an operation destroys the integrity of the panel edge, and thus renders the panel essentially useless for edge compression strength testing, and other test measuring purposes.\nThere is a clear and strong need for a corrugated paperboard panel cutting method and apparatus, wherein the first cut of the panel test specimen, is made either precisely parallel, or precisely normal, to the flute direction, such that the final square, or other rectangular, test speimen panel, has two edges running precisely normal to the flute direction and two edges running precisely parallel to the flute direction. If the test specimen's first cut edge can have the desired precise orientation relative to the flute direction, the remaining three test specimen edges can automatically be made to have the necessary orientations, e.g., by using T-square cutting guide surfaces properly oriented to the cutter blade plane. The present invention is directed to a method and apparatus for achieving either the desired parallelism, or normalcy, between the test specimen's first cut edge and the flute direction, as well as subsequent cut edges.\nAn additional problem inherent in conventional cutter blade constructions, which has been elegantly solved by the present invention, is the premature dulling of the cutter blade. With conventional arrangements, the same zone of the cutter blade's cutting edge is used continually for each cutting cycle. The cutting pressure is repetitively applied to a relatively small localized area of the cutter blade edge, so that the cutter blade becomes dull after only a relatively few cutting cycles, necessitating frequency replacement of the cutter blade.\nCutter blade replacement requires time, and also involves potential danger to the operator, because the cutter blade edge is extremely sharp, especially the new replacement cutter blade. Manual handling of the cutter blade during blade removal or installation can, therefore, cause serious injury to the operator.\nThe present invention also seeks, in part, to provide a cutter blade construction, wherein, different zones of the cutter blade's cutting edge are used for cutting purposes during successive cutting cycles. Wear on the cutter blade is, therefore, distributed over a longer cutting edge area, such that cutter blade replacement is less frequent, compared to prior art arrangements."}
{"text": "In certain system arrangements duplicated buffer memory units are employed in an active-standby configuration to reduce so-called \"down time\" in the presence of a failure and, thereby, improve system reliability. In such duplicated memory system arrangements proper system operation requires that the contents of each of the memory units be identical. Prior arrangements are known in which synchronization of such duplicated memory units is realized. (See, for example, U.S. patent application Ser. No. 07/739,928, filed Aug. 2, 1991 now U.S. Pat No. 5,278,969 issued Jan. 11, 1994.) It should be noted that ATM transmission of data is on a cell-by-cell basis wherein each cell comprises a fixed length packet. Thus, in the ATM context, the terms cell and packet are interchangeable. In the noted prior duplicated memory arrangement, a plurality of queues is employed to store cells, with each of the queues being associated on a one-to-one basis with either inputs or outputs of the memory unit. Synchronization of the active and standby memory units is realized in this prior arrangement by employing a queue length counter with each queue and observing when all of the counters have reached zero count at least once. Although this prior arrangement will indicate that the duplicated memories are in synchronization at some particular time after synchronization has been initiated, it does not determine when or how to initiate resynchronization. Once synchronization has been obtained, if the contents of the corresponding queues in the active and standby memories become different during any cell interval, problems can arise."}
{"text": "1. Field of the Invention\nThe present invention relates, in general, to a magnetoresistive (MR) read/write head and method for making a MR read/write head, and, more particularly, to a MR read/write head having decreased track width and method for making a MR read/write head with high control over track width definition.\n2. Relevant Background\nThe desktop personal computer market continues to demand higher capacity and faster performance from hard disk and tape drives. With applications such as file downloading, increased file sizes, advanced operating systems, and multimedia applications, demand for hard disk drive capacity, for example, is doubling every year. Technologies for storing and retrieving data from magnetic media must also be cost effective. Because lower cost per megabyte (MB) is also desired, the prior practice of simply adding more disks and \"heads\" (i.e., structure in which read and write elements are provided) to a hard drive is less and less effective. Disk and tape drive suppliers continue to increase areal densities, or the number of data bits per square inch, to meet the increasing demand for storage at competitive pricing. Read and write head design are key technologies needed to achieve these capacity increases.\nThe write element that writes data on the disk is typically made up of two poles that are separated by a write gap, and which generate a magnetic field when they are excited by a coil magnetically coupled to the poles. When the write element is in proximity to the disk, a magnetic field generated by the poles sets the magnetic orientation in given locations on the disk. In this manner, data is written on the disk.\nThe read element that reads data from the disk is sandwiched between two shields. During a read operation the read element flies in proximity to the disk so that the read element senses the magnetic orientation of the given disk locations. To enable the read element to focus on a small disk location during reading at (i.e. the read element must not be affected by the magnetic orientation of adjoining disk locations), it is desirable to shield the read element. The two relatively large shields filter out the magnetic effects of adjoining disk locations, so that a specific disk location can be focused upon for reading.\nHard disk drives with lower areal densities typically use inductive read and write elements. Inductive heads offer low cost and mature processing technology suitable for high volume production. To increase the signal strength from an inductive head, designers have increased the number of turns in the read head as the read signal is directly proportional to the number of turns. Some inductive heads use fifty or more turns in the read/write head. However, increasing the number of turns increases the head's inductance. There is a limit to the amount of inductance a head can tolerate to effectively perform data write operations. Since thin-film inductive heads use the same inductive element for both reading and writing, the head cannot be optimized for either operation. Moreover, the increased inductance decreases the frequency at which data can be written to and read from the magnetic media.\nMagnetoresistive (MR) head technology is used to provide higher areal density than possible with inductive heads in both disk and tape drives. MR head structures include an MR element as a magnetic field sensor. A coil is formed above the read head and magnetically coupled to the magnetic yoke that defines the poles of the write element. Although the coil and yoke are magnetically coupled, they are separated by an insulating material to prevent current flow between the coils and the yoke. To provide an area efficient structure, it is desirable to vertically stack the coils in two or more layers.\nAn MR head generally combines the read and write elements of the head into an integrated unit. It does so by eliminating one of the poles of the write element and substituting in its place one of the shields of the read element. In doing so an integrated pole/shield element is created.\nUsing an MR structure as a read element provides high signal output and low noise compared to inductive heads. This higher signal output allows the write element to write data in a much narrower track while still being reliably detected by the MR read element. Separate read and write heads allow each head to be optimized for one particular function (i.e., reading or writing data). With an MR head, the number of wire turns in the write element can be greatly reduced, resulting in a low inductance head enabling high frequency write operations.\nThe track width of an MR head is largely determined by the size of the area of the disk that is affected by the write head. Where the pole/shield structure is physically large, the pole/shield will tend to undesirably affect a larger part of the disk during a write operation, which is a phenomenon referred to as \"fringing\". Fringing has an adverse effect on the efficient storage of data on the disk given that it is usually desirable to pack data on the disk as densely as possible, thereby increasing the storage capacity of the disk.\nThe track width can be decreased by making the poles physically small at the write tip (i.e., the portion of the yoke that forms the poles), thereby concentrating the magnetic field into a smaller area. However, in conventional MR head processes, the yoke, including the write tip, are formed as an integrated structure over the coil structure. The coil structure is very thick, especially when vertically stacked coils are used. Hence, the write tip is typically patterned using thick photoresist (on the order of 10-15 microns thick) making it difficult to define the small structures that are required to decrease track width. Critical dimension control is poor when patterning thick layers of photoresist resulting in unacceptable variation in the size of the patterned feature.\nWhat is needed therefore, is an MR head that combines the advantages of a small write tip structure, but that can be manufactured with a high degree of process control."}
{"text": "1. Field of Invention\nThe invention relates to a cooling device and, in particular, to an air flow channel.\n2. Description of the Prior Art\nFIG. 1 shows the cooling method for a conventional projector lamp. Normally, a blower 800 directly blows cooling air into the lamp reflector 700, taking away heat generated by the projector lamp 500. The cooling places of the projector lamp 500 are mainly the high-temperature points above and below the filament 510, the soldering point 520, and the end points 530. Each high-temperature point has to be controlled within a certain temperature range. Aside from the projector lamp 500, the lamp reflector 700 has to be cooled below a predetermined temperature in order to prevent the lamp reflector from breaking. Using the conventional cooling method on the projector lamp 500, the cooling air has a directly impact on the projector lamp 500 instead of cooling each of the high-temperature points according to the required temperature ranges. Moreover, the cooling air first flows by the end point 530, the soldering point 520, and the filament 510 and absorbs their heat before reaching the bottom of the lamp reflector 700. However, these high-temperature points have temperatures higher than the bottom of the lamp reflector 700. The cooling air after absorbing heat from them is very difficult to further cool down the bottom of the lamp reflector 700. If one solves the problem by increasing the fan speed, the temperatures are those high-temperature points will be lowered too much out of the required ranges. This method of first cooling down the high-temperature points with higher temperatures before cooling down the lamp reflector has the drawback that when the temperature of the lamp reflector becomes too high, adjusting the fan speed to lower the temperature of the lamp reflector will render the temperature of the filament out of the specific range. This makes the temperature control very difficult.\nAs described, the prior art results in non-uniform cooling on different parts of the projector lamp 500 and the lamp reflector 700 or cannot cool the system in such a way that the projector lamp 500 and the lamp reflector 700 fall right within the desired temperature ranges. These effects will affect the performance and lifetime of the projector lamp and the lamp reflector. Thus, it is desirable to provide a better cooling system."}
{"text": "1. Field of the Invention\nThis invention primarily relates to a guiding means for an endless apron provided in the drafting part of a spinning frame, and further, to a guiding means for an endless belt in the feeding mechanism of, for example, a printing machine, duplicator, precision packaging machine, film processor, and the like.\n2. Description of the Prior Art\nFIG. 1 is an illustration showing a typical apron zone of a long apron type drafting device which comprises such components as a bottom roller BR and a top roller TR, both being disposed on a roller stand RS which supports a roller, a cradle CR, a tensor bar TB, a tension guide TG pressed by a tension lever TL in the direction of the arrow, a top apron TA and a bottom apron BA, in addition to a weighting arm WA.\nHowever, the above said apron guiding device in respect of performance and operability for tightening, balancing and guiding the bottom apron, as well as performance of the tensor bar itself and operability in maintenance thereof, has not proved satisfactory."}
{"text": "1. Field of the Invention\nThe present invention generally relates to the manufacture of electrical conductor lines and vias that interconnect circuits on substrates such as semiconductors and related packages and, more particularly, to a low cost method of filling seams or holes in substrates using a combination of a low resistivity metal deposited by physical vapor deposition (PVD) and a refractory metal deposited by chemical vapor deposition (CVD). The invention has particular application in submicron circuit manufacture.\n2. Description of the Prior Art\nLow resistivity metals such as aluminum and copper and their binary and ternary alloys have been widely explored as fine line interconnects in semiconductor manufacturing. Typical examples of fine line interconnect metals include Al.sub.x Cu.sub.y, where the sum of x and y is equal to one and both x and y are greater than or equal to zero and less than or equal to one, ternary alloys such as Al--Pd--Cu and Al--Pd--Nb, Al--Cu--Si, and other similar low resistivity metal based alloys. Today's emphasis on scaling down line width dimensions in very large scale integrated (VLSI) circuitry manufacture has led to reliability problems including inadequate isolation, electromigration, and planarization.\nThe IBM Technical Disclosure Bulletin to Ahn et al., Vol. 33, No. 5, October 1990, pages 217-218, discloses tungsten wrapped copper conductors and via holes fabricated by selective deposition using a mixture of WF.sub.6 and SiH.sub.4 in the presence of hydrogen. Encapsulated interconnects like those of Ahn et al. have significantly higher resistance to electromigration and the small grain size of a selective tungsten film reduces reflectance and thereby enhances the ability of photolithography tools to focus and resolve photoresist images. However, the tungsten layer formed using the low temperatures described by Ahn et al. would be silicon rich (e.g., 3-4%) and would not be a good diffusion barrier for copper since copper resistivity would be degraded by the formation of copper silicide. Thus, it is difficult to deposit a diffusion barrier by selective means at low temperature. Moreover, the Ahn et al. technique relies on the formation of a donut shape at the bottom of the lines which is normally created by the reaction of outgassing moisture and WF.sub.6. The creation of the donut shape is believed to be not reliable.\nDalton et al., VMIC Conference, Jun. 12-13, 1990, pages 289-292, points out that a hot wall CVD reaction involving SiH.sub.4 and H.sub.2 reduction of WF.sub.6 to form a selective tungsten layer on an aluminum or alloy conductor results in the incorporation of fluorine at the aluminum and tungsten interface. The fluorine incorporation is a byproduct of the reaction of WF.sub.6 with aluminum as shown by Equation 1. EQU WF.sub.6 +2Al.fwdarw.2AlF.sub.3 +W Eq. 1\nThe thin layer of aluminum fluoride will increase the series contact resistance of Metal 1 to Metal 2 vias. Dalton reported that sputtering TiW film on top of the aluminum prior to tungsten encapsulation using CVD eliminates the problem of fluorine absorption.\nDalton discloses a traditional scheme for interconnect formation wherein aluminum is first deposited on a planar surface, it is overcoated with the sputtered TiW layer (the only difference from traditional processing), the aluminum is then patterned using photoresist imaging and developing followed by reactive ion etching (RIE). The resulting structure is then overcoated with a passivation dielectric such as SiO.sub.2 or polyimide which itself is subsequently patterned, subjected to RIE, and metallized to create a multilayered structure. FIG. 1 is taken from Dalton and shows that multilayer devices produced by traditional processing schemes have seams in the dielectric layers at the location of the metal conductor lines and have a very irregular top surface.\nIt is difficult to achieve planarity of the dielectric using RIE. Planarity is in part dependent on the pattern density, and non-planar surfaces result in puddling problems during subsequent metalization. If an RIE technique is used on polyimide, an etch stop is needed for removal of photoresist on top of aluminum or copper based lines when the lines are etched down to the polyimide surface because the photoresist removal process would also remove polyimide. RIE of any high copper content, aluminum or copper alloy is extremely difficult. A serious drawback of traditional processes which include metal RIE is that a large number of metal shorts tend to develop with fine geometry due to particle defects.\nU.S. Pat. No. 4,824,802 to Brown et al. discloses a method for filling interlevel dielectric vias or contact holes in multilevel VLSI metalization structures. In particular, an intermediary metal such as tungsten or molybdenum is either selectively deposited in openings in an insulator or non-selectively deposited over the entire surface and in the openings of the insulator by CVD, then a planarization resist, such as azoquinonenovolac-type resists, polymethacrylates, polyimides, or other thermoplastic materials, is applied over the top of the intermediary metal. A planarized structure is then obtained by etching to a level where the intermediary metal is even with the resist. The Brown et al. method does not avoid metal corrosion and other problems associated with etching and is not useful for planarizing Al--Cu or other soft alloys because they have different properties from the harder metals such as tungsten and molybdenum. Moreover, using the Brown et al. method, it is difficult to completely fill vias and lines.\nU.S. Pat. No. 4,944,836 to Beyer et al. discloses a chemical-mechanical polishing technique which can be used to produce coplanar metal/insulator films on a substrate. In particular, Beyer et al. contemplate patterning an underlying insulating layer, depositing an Al--Cu film, and then using a chemical-mechanical polishing technique wherein an alumina slurry in dilute nitric acid is mechanically rubbed on the surface to remove Al--Cu. The polishing compound tends to have a significantly higher removal rate for Al--Cu than the underlying insulator. The resulting structure includes Al--Cu lines planarized with the insulating layer, and subsequent layers can easily be added in the fabrication of multilayer structures.\nU.S. Pat. No. 4,956,313 to Cote et al. discloses a via filling and planarization technique wherein Al--Cu alloy lines are patterned on top of a first passivation layer on a substrate, the lines are overcoated with a second passivation layer which is preferably a doped glass such as phosphosilicate glass (PSG) or borophosphosilicate glass (BPSG) which conforms over the contours of the Al--Cu alloy lines, vias are then formed in the second passivation layer to expose the lines, and tungsten is applied over the surface of the second passivation layer and in the vias by CVD. It is reported in Cote et al. that CVD tungsten is conformal in character and can fill the vias without creating voids. The structure is then planarized by polishing with an abrasive slurry.\nNeither Beyer et al. nor Cote et al. recognize that polishing is not practical for low resistivity, soft metals such as Al--Cu alloys. This is because such materials tend to scratch, smear on the surface, and corrode under the influence of the slurry. Moreover, creation of the planarized structures in accordance with Cote et al. takes several processing steps which increases costs and reduces output.\nRossnagel et al. J. Vac. Sci. Technol. 2:261 (Mar/Apr. 1991) discloses a collimated magnetron sputter deposition technique for depositing films that are compatible with lift-off patterning techniques and hole filling. The technique is also presented in U.S. Pat. No. 4,824,544 which is herein incorporated by reference.\nShiozaki et al., Abstracts of the 19th Conference on Solid State Devices and Materials, discloses the use of selective tungsten deposition for hole filling on top of a high resistivity hard metal such as MoSi.sub.x and is unrelated to encapsulation of a soft metal."}
{"text": "This invention was made in part with support from the National Science Foundation Grant Number DMR 9315914.\nThe present invention relates to a process for forming high quality crystalline refractory materials, particularly gallium (Ill) nitride (GaN), from solid precursors. GaN is a material newly available for use in the optoelectronics industry fur the fabrication of light-emitting diodes (LEDs) and blue lasers. It is also possible that doped GaN crystals may have utility as semiconductors. A particularly suitable application is the replacement of standard light bulbs in large outdoor displays, traffic lights and street lighting by GaN LEDs. GaN crystals, when properly activated, fluoresce producing a bright blue glow which is about 60 times brighter than the best GaP based yellow-green LEDs and many times brighter than a standard light bulb which it would replace. Further, a GaN LED display would have an operating life far in excess of the standard light bulb.\nCurrently, bulk quantities of high purity, polycrystalline gallium nitride are not available. Current techniques to produce such materials require maintaining reactants at high temperatures and pressures for long periods of time. Prior attempts to manufacture GaN by reacting gallium iodide with lithium nitride, which appears to be a suitable approach, produces elemental Ga, nitrogen and Lil and not GaN. Thus there is a need for a low cost, rapid process to produce large quantities of powdered crystalline materials for use in such applications as lighting, signal displays, and flat screen displays for computers and television screens."}
{"text": "Hyaluronic acid is a naturally occurring linear polysaccharide. It is a polymer of disaccharide units composed of β-1,4-D-glucuronic acid and N-acetyl-D-glucosamine linked together by β-1,3-glycosidic bonds. Hyaluronic acid has a molecular weight ranging from 1,000 to 10,000,000 daltons.\nHyaluronic acid is the only non-sulfated glucosaminoglycan (GAG) found in the extracellular matrix (ECM) of higher animals.\nBiomaterials derived from hyaluronic acid are used, for example, in regeneration of cartilage and skin. Also, they play key roles in drug delivery or surgery. Chemically unmodified hyaluronic acid is utilized as the aid for the delivery of ophthalmologic drug to improve the absorption of drugs and proteins at the mucosal tissue.\nFurther, with regard to the treatment of osteoarthritis, hyaluronic acid improves lubrication at the joint surface and thus reduces pain. Arthritis reduces the production of hyaluronic acid and the increased breakdown by proteases further reduces hyaluronic acid in the joint. As a result, the articular damage can be aggravated because external impact cannot be adequately absorbed or distributed at the joint.\nSince approved by the FDA in 1997, the injection of hyaluronic acid is widely used. With compositions similar to those of the joint fluid, hyaluronic acid solution regains viscosity and elasticity when injected into the joint and improves lubrication and absorbs impact, thereby protecting the joint cartilage and preventing further damage.\nAlso, it is reported to have anti-inflammatory effect and inhibit the breakdown of chondrocytes by several in vitro researches. Based on this, it can be inferred that hyaluronic acid can be utilized as scaffold for cell delivery in the process of cell transplantation therapy of chondrocytes for cartilage tissue regeneration.\nThe in vivo metabolism of hyaluronic acid is regulated by hyaluronic acid synthase (HA synthase) and hyaluronidase. Hyaluronic acid synthase is an integral membrane enzyme, which synthesizes hyaluronic acid polymers and excretes them out of the cell. In mammals, HAS1, HAS2 and HAS3 isoenzymes are found.\nHyaluronidase is an enzyme that breaks down hyaluronic acid. Depending on the activation pH, it can be classified into neutral hyaluronidase and acidic hyaluronidase. Neutral hyaluronidase (PH-20) is specifically found in the testis (sperm) and shows activity in the physiological pH.\nAcidic hyaluronidase (Hyal1-4) is found in body fluids as well as in a variety of human organs, including spleen, cartilage, skin, eyes, liver, kidney, bladder, placenta, etc., and shows activity around at pH 3. Since this enzyme is present in the lysosome, hyaluronic acid is transported into the cytoplasm by endocytosis and degraded.\nThe endocytosis of hyaluronic acid is mediated by the cell surface receptors. The primary cell surface receptor for HA is CD44. CD44 is a transmembrane glycoprotein and binds primarily to hyaluronic acid. It is expressed in most of human cell membranes. The hyaluronic acid fragments degraded at the lysosome induce different signal transfers depending on their size, affecting cell proliferation and differentiation.\nVarious porous scaffolds are used in the field of tissue engineering. Chondrocytes cultured in 2D condition such as petridish cannot maintain their chondrocytic phenotype, resulting in synthesis of type I collagen rather than type II.\nIn contrast, chondrocytes cultured in 3D condition maintain their characters and functions, resulting in spherical cell shape. That is, a 3-dimensional scaffold helps the cell attachment and cartilage formation in vitro and in vivo.\nThe use of porous scaffolds is very useful in that chondrocytes can be cultured under a restricted condition with significantly reduced necrosis of the tissue.\nThus, as one of new treatments for joint cartilage damage, the technique of culturing stem or progenitor cells derived from the cartilage, bone marrow or periosteum in vitro using biodegradable scaffolds and transplanting them into lesion site is attempted instead of direct injection of autologous chondrocytes.\nThis technique is advantageous over the method of using cell suspension since necrosis of the cartilage tissue to serve as cell sources is significantly reduced and the cells can be incubated under a restricted condition. Further, it is also adequate considering that cell attachement is improved by a scaffold.\nIn addition to cell delivery, a scaffold can serve as mold to fill lesion site of tissue. An ideal scaffold is required to be non-immunogenic, non-toxic, biocompatible, biodegradable and easy to manipulate.\nThe pore size of the scaffold is important. An inadequate pore size may result in the restricted nutrient supply to the cells being cultured in the scaffold or ineffective removal of the wastes from the cells.\nIf the pore size is too small, the cells may not be completely filled inside the scaffold, which negatively affects the generation of the tissue. In contrast, if the pore size is too large, not all the extracellular matrix proteins synthesized by cells are accumulated. That is, the proteoglycan and collagen accumulated in the scaffold can be found only a part of the total degree of actually synthesized ECM.\nAccordingly, even though chondrocytes cultured in a porous scaffold in a metabolically active state, the scaffold may deem as inadequate for the growth of chondrocytes.\nThus, the control of the pore size is important in improving the retention rate of newly synthesized extracellular matrix molecules inside the scaffold.\nAlso, in case the sponge type scaffold with minimal porosity required for survival, growth and differentiation of cells is in solid state, it comes into problem that the surgical operation is the only way of delivering the entire scaffold including cells to a disease site."}
{"text": "Recently, IP-related protocol groups such as not only IP (Internet Protocol) but also SIP (Session Initiation Protocol; RFC3261) to execute call control among radio communication terminals have been used also in a radio communication system such as a mobile communication system.\nIn a radio communication system using the protocol groups, a group call is provided in which a call within a group consisting of multiple radio communication terminals is established when one of the radio communication terminals (for example, a cellular telephone terminal) calls up the other radio communication terminals. Each of the radio communication terminals participating in the group call can transmit and receive voice data and video data to and from the other radio communication terminals participating in the group call.\nMeanwhile, there is proposed a radio communication system which allows a radio communication terminal to participate halfway in a group call (hereinafter such participation being referred to as halfway participation) (see Patent Document 1). Specifically, in this radio communication system, a control station device periodically broadcasts notification information including information on the group call to radio communication terminals through a radio base station.\nIn this way, a radio communication terminal which has failed to receive a call-up signal can receive the notification information, and participate halfway in the group call by transmitting a request for participation in the group call. Patent Document 1: JP-A 2001-168793 (pages 5 to 7, FIGS. 1 and 3)"}
{"text": "The purpose of the invention is to saw particularly hard stones with two parallel blades tensioned in a frame having an almost vertical oscillating motion.\nFor hard stone sawing like granites, one commonly uses sawing machines with one or more blades, animated by a straight reciprocal or oscillating movement. The blades used are either steel ones, sprinkled with water and abrasive grits (like sand, steel shot or silicon carbide) either ones with diamond segments.\nSome of those machines only have one blade animated by a horizontal axis oscillation while the stone blocks are in horizontal traverse movement. The blade is tensioned in a frame to avoid deviations during the sawing. It is better to place the blade laterally in order to accomodate blocks of unlimited width.\nThe major inconvenience of such machines is that one has to cut one slab at a time, rendering the sawing slow and expensive. Another inconvenience is given by the asymmetrical disposition of the blade in tne frame, introducing a bad equilibrium of the tensions in the machine and this creates deviations during the sawing.\nMachines having several blades are also well known. Those blades are tensioned parallely in a frame which is animated by a horizontal axis oscillation or a straight reciprocal vertical movement, with a horizontal traverse movement for the block to be sawn. The distance between the blades is chosen in order to saw the slabs with a desired thickness. The evident advantage of such machine is that series of slabs can be sawn in one operation.\nThe inconveniences of such machine are unfortunately numerous. The tensioning of the blades in the frame is, indeed, a long and arduous operation and it is not possible to modify the distance between the blades as one would like. The width of the sawn slabs can not be regulated in accordance with clients' demands and one has to be satisfied with standard width.\nAnother inconvenience of multiblade machines lies in the frame wherein the block to be sawn has to pass through. One can thus see that the frame dimensions are determined for those of the blocks. In particular, the width of the blocks has always to be below the inside width of the frame.\nIt must also be observed that the construction of such kind of multiblade machine is sophisticated and always expensive.\nThe object of the hereby invention is to remedy the above mentioned inconveniences, realizing a double blade sawing machine, giving the opportunity of simultaneous sawing of two slabs with arbitrary thickness in two independent blocks placed on both sides of the machine. One realizes like this a symmetrical machine, fast and economical."}
{"text": "There is a need for low bend loss optical fibers, particularly for optical fibers utilized in so-called “access” and fiber to the premises (FTTx) optical networks. Optical fiber can be deployed in such networks in a manner which induces bend losses in optical signals. transmitted through the optical fiber. Some applications that can impose physical demands, such as tight bend radii, compression of optical fiber, etc., that induce bend losses include the deployment of optical fiber in optical drop cable assemblies, distribution cables with Factory Installed Termination Systems (FITS) and slack loops, small bend radius multiports located in cabinets that connect feeder and distribution cables, and jumpers in Network Access Points between distribution and drop cables. It has been difficult in some optical fiber designs to achieve both low bend loss and low cable cutoff wavelength at the same time."}
{"text": "1. Field of the Invention\nThe present invention pertains to a pressure plate assembly for a friction clutch, which assembly includes a housing; a pressure plate connected to the housing for rotation in common around an axis of rotation which is free to move axially with respect to the housing; a stored-energy device which pretensions the pressure plate in the axial direction with respect to the housing; a wear-compensating device which acts in the path of support between the stored-energy device and at least one of the housing and the pressure plate and comprises at least one wear-compensating element, which is free to move when wear occurs and/or when wear is compensated; an indicator arrangement comprising at least one indicator opening in the housing and at least one indicator element on a wear-compensating element, where information concerning the amount of wear which has occurred can be derived from the position of the indicator element relative to the minimum of one indicator opening.\n2. Description of the Related Art\nU.S. Pat. No. 5,531,308 discloses a friction clutch in which a wear-compensating device can be used to ensure that the wear which occurs during the operation of the clutch, such as the wear of the friction linings of a clutch disk, for example, is compensated. For this purpose, two adjusting rings are provided, one of which is held nonrotatably with respect to a housing, whereas the other adjusting ring is free to rotate to compensate for wear. The two rings have corresponding ramp surfaces, so that, upon rotation of the second adjusting ring, the total axial length of the wear-compensating device consisting essentially of the two adjusting rings changes. On the adjusting ring which can rotate relative to the housing, an axial projection is provided, which engages in a curved, slot-like opening in the housing, i.e., in a bottom area of the housing. When the adjusting ring turns, this projection moves in the opening through which it passes. From the position which this projection occupies in the opening, it can be determined how far the rotatable adjusting ring has rotated up to that point, which is an indication in turn of the extent to which wear compensation has occurred and thus also of the amount of wear which has occurred.\nIt is the task of the present invention to improve a pressure plate assembly of the general type in question in such a way that, by the use of a simple design, the amount of wear which has occurred can be detected more easily.\nAccording to the invention, the indicator element at least partially fills at least one indicator opening in at least one state of wear.\nWith the design of a pressure plate assembly according to the invention, the geometry of the indicator element can be made essentially independent of the geometry of the indicator opening. This facilitates the design and also makes it possible for both the indicator element and the indicator opening to be designed with respect to their shape and/or configuration in such a way that, through the cooperation of these two assemblies, a very precise indication of the actual amount of wear which has occurred or of the degree of wear compensation which has occurred can be obtained.\nIt is possible, for example, for the width of the indicator element, i.e., the dimension essentially transverse to the direction of its movement with respect to the minimum of one indicator opening, to change over the length of the indicator element in the direction of movement. It is thus achieved that the width of the indicator element visible in the minimum of one indicator opening changes as the element moves past, so that, as the indicator element moves, the area of the indicator opening which it fills also changes, from which information can be derived concerning different states of wear.\nFor this purpose, it is possible, for example, for the width of the indicator element to change continuously, at least in certain of its areas. Alternatively or in addition, it is possible for the width of the indicator element to change in step-like increments. Although a step-like change in width produces a very clear-cut discontinuity in the indicator characteristic, an indicator element with a continuous change in width provides information of higher resolution concerning the state of wear.\nIt is also possible to provide a plurality of indicator openings in succession in the direction of the movement of the indicator element. In another, preferred embodiment of the pressure plate assembly according to the invention, it is possible, in conjunction with the provision of a plurality of indicator openings, for the indicator element to have an indicator section which fills essentially all of the indicator openings in one wear state and at least partially fills one indicator opening or essentially no indicator opening in at least one other wear state. It can therefore be provided, for example, that the indicator fills all of the indicator openings when the clutch is new and has not yet been affected by wear, and that, as the amount of wear increases, the indicator element moves successively along the indicator openings and thus opens various indicator openings one after the other.\nAs an alternative, it is also possible for the indicator element to have an indicator section which does not fill all of the indicator openings in any wear state. It can be provided in this case, for example, that, as a function of the amount of wear which has occurred, the indicator section completely fills one of the indicator openings or partially fills two of the indicator openings.\nIn a design of the pressure plate assembly which is very simple to realize, it is possible for the wear-compensating device to comprise an adjusting ring, which can rotate around the axis of rotation for the purpose of wear compensation, and for the indicator element to be free to move along with the adjusting ring constituting the wear-compensating element.\nAn alternative embodiment of the pressure plate assembly according to the invention can be designed in such a way that the wear-compensating device comprises: an arresting element which, upon the occurrence of wear, is able to move with respect to the pressure plate assembly in correspondence with the amount of wear which has occurred; an essentially wedge-shaped slider element, which is installed between the pressure plate assembly and the arresting element and which, upon movement of the arresting element with respect to the pressure plate assembly, is free to move under the pretensioning effect in correspondence with the extent of the movement of the arresting element with respect to the pressure plate assembly in order to keep the arresting element in the position relative to the pressure plate assembly to which it has moved as a result of wear; and an adjusting ring, which is able to rotate around the axis of rotation to perform a wear-compensating movement, where the adjusting rotation of the adjusting ring is limited by the slider element and/or by the arresting element, where the indicator element is able to move along with the slider element forming the wear-compensating element.\nIn addition, the pressure plate assembly according to the invention can be designed so that the stored-energy device is located in an area between the indicator element and the minimum of one indicator opening and has a see-through opening associated with the minimum of one indicator opening. To avoid the introduction of additional openings in the stored-energy device, it is possible, for example, for the see-through opening to be formed by an area between two of the spring tongues of the stored-energy device.\nIt is also possible, according to the invention, for the indicator element to be connected to the wear-compensating element by connecting sections provided on the indicator element. It is also possible for the indicator element to have a contact area, by means of which it is supported on a surface of the wear-compensating element facing the axis of rotation. It is possible in this way to prevent the indicator element from shifting position as a result of centrifugal force during rotational operation.\nThe present invention also pertains to a friction clutch which contains a pressure plate assembly according to the invention.\nOther objects and features of the present invention will become apparent from the following detailed description considered in conjunction with the accompanying drawings. It is to be understood, however, that the drawings are designed solely for purposes of illustration and not as a definition of the limits of the invention, for which reference should be made to the appended claims. It should be further understood that the drawings are not necessarily drawn to scale and that, unless otherwise indicated, they are merely intended to conceptually illustrate the structures and procedures described herein."}
{"text": "(a) Fields of the Invention\nThe present invention relates to solid-state image sensing devices having a solid-state image sensing element and a transparent substrate for protecting the element, and to their fabrication methods.\n(b) Description of Related Art\nIn a solid-state image sensing device employing a CCD (Charge Coupled Device) and the like, the area of a photodiode serving as a light receiving unit has been decreasing by demands for downsizing and resolution enhancement thereof. Such a decrease in the area of the light receiving unit in turn degrades the light collection efficiency of the device. In order to make up for this degradation, a so-called microlens has come to be developed and then put into use. This microlens is typically made of resin and disposed above the light receiving unit formed on each pixel. The microlens refracts light not coming directly in the light receiving unit to collect the refracted light into the light receiving unit, thereby enhancing light collection efficiency to improve sensitivity.\nFIGS. 15 and 16 show cross-sectional structures of a conventional solid-state image sensing device. Referring to FIG. 15, in the surface of a substrate 101 for a CCD-type solid-state image sensing element, a recess is provided on each pixel. At the bottom of the recess, a photodiode 102 is provided which converts an incoming light into an electrical signal. On the substrate 101 for the solid-state image sensing element, a first acrylic flattening film 103 is formed which flattens unevenness of the substrate surface. On the first acrylic flattening film 103, color filters 104 are formed to be associated with the photodiodes 102, respectively. On the color filters 104, a second acrylic flattening film 105 is formed which flattens unevenness generated due to gaps between the color filters 104. On the second acrylic flattening film 105, microlenses 106 are formed to be associated with the photodiodes 102, respectively.\nReferring to FIG. 16, a solid-state image sensing element 113 is composed of the photodiodes 102, the color filters 104, the microlenses 106, and the like formed on the substrate 101 for the solid-state image sensing element. The solid-state image sensing element 113 is mounted inside a package 112, and the top of the package 112 is covered with a transparent substrate 109. As shown in FIG. 15, inside the package 112, an air space 110 is interposed between the solid-state image sensing element 113 and the back surface of the transparent substrate 109. As shown in FIG. 15, when light 111 transmits through the transparent substrate 109 to come in the microlenses 106, reflections occur on the top and bottom surfaces of the transparent substrate 109 and the top surfaces of the microlenses 106.\nHowever, rapid downsizing of the device in recent years has made it difficult to secure sufficient sensitivity only by light collection from the microlenses. To overcome such a difficulty, a structure such that an antireflection film is formed on the microlenses is proposed (see Patent Document 1: Japanese Patent No. 2719238).\n(Patent Document 2)\nJapanese Examined Patent Publication No. 7-54974\n(Patent Document 3)\nJapanese Examined Patent Publication No. 7-28014\n(Patent Document 4)\nJapanese Patent No. 2942369"}
{"text": "Concrete is traditionally used for floors in both residential and commercial applications in view of its robustness and economic benefits. Depending upon the circumstances, the concrete may be left unfinished, partially finished or completely finished wherein a high gloss decorative service is obtained.\nIn warehouses, factories, etc., concrete floors are routinely cleaned by rotary driven machines that employ brushes located on the underside of the machinery whereby the machinery traverses the floor to provide a clean surface. Typically these floor cleaning machines have a tendency of progressively deteriorating the surface of the concrete floor. This is caused by the bristles extending into the naturally occurring crevices in the concrete floor thus causing minute particles of the concrete to break away. Through the repeated cleaning utilizing this floor polishing machine is then typically employed to treat the surface of the floor so as to restore the floor back to a desirable appearance.\nTypically the flooring machines that are used on a routine basis sit idle as they have a limited purpose, which is to clean the floor. These traditional flooring machines are not used to polish the floor, but merely for the single purpose of cleaning the floor. Likewise, a separate floor polishing machine traditionally is used for the sole purpose of polishing the concrete floor when the floor condition has sufficiently deteriorated. Accordingly, because these machines have limited purposes, they are seldom used, thus creating inefficiencies for a business.\nIt will be desirable to forego the aforementioned costly steps by providing an improved cleaning and honing brush that works as an attachment to an ordinary cleaning machine, such as a Tennant or Advance brand scrubber machine. It would be desirable to provide an improved cleaning and honing brush that operates under low pressure, does not require an independent power source, continuously exposes new abrasive material during the cleaning process and has interchangeable replaceable polymer brush strips that can be easily removed and replaced with replacement polymer brush strips.\nThe aforementioned problems may be overcome by providing a polymer brush strip that may be co-extruded with diamond particles integral with the head of the brush strip.\nIt will further be desirable to provide a floor resurfacing device that operates under low speed and low pressure conditions while utilizing a diamond impregnated brush made of polymer matrix that can be used on ordinary rotary machines including low power, low pressure automatic flooring machines.\nIt will be desirable to provide an improved polishing system that can be used in connection with a traditional automatic flooring machine wherein a series of diamond embedded polymer strips are used with a rotary disc. The rotary disc bearing the polymer strips can also be used for cleaning the surface of the concrete as well as providing a polishing aspect during the process of cleaning. The floor cleaning process employs a rotary disc having polymer strips with greater grit density so as to provide an improved surface finish. The process is repeated routinely whereby at each cycle the polymer strips are changed out to have finer grit size so as to continuously enhance the surface finish quality. The aforementioned process is accomplished while cleaning the floor during low speed and low pressure conditions. Said process can be used with floors other than concrete floors.\nThe present invention also solves a problem that is found in the maintenance of floors in offices. For example, in office settings a common floor covering is vinyl composite tile (VCT). Wax is often applied over the tile so as to protect its surface and to provide an enhanced floor appearance. Over time these floors acquire a build up of wax and other particles that need to be removed so as to reestablish the floor's appearance to a new-like condition. Traditionally, a common way of stripping wax and other particles from floors was to use chemicals. The chemicals loosen the wax so it can then be removed via a scrubbing process. Typical ways of removing the loosened wax were to use an abrasive pad. However, the loosened wax tends to gum up the surface of the pad making removal very difficult. It would be preferred to provide a more environmentally green method of stripping wax from VCT, for example, use of water would be preferred. In order to accomplish complete removal of wax from VCT, an improved pad driver with brush design would be desirable. It would also be desirable to provide an improved brush that has a polymer component with a metal component."}
{"text": "1. Field of the Invention\nThis invention relates to a three way sliding poppet valve for switching high pressure, high temperature fluids under low pressure control and, in particular, this invention relates to a low internal volume three way valve for high speed switching of high pressure, high temperature fluids without volume change for precision laboratory measurements such as is encountered in permeameter tests with corrosive fluids involving core samples.\n2. Description of the Prior Art\nIn permeameter tests, high temperature (such as 280 degrees Fahrenheit) and high pressure (such as 10,000 psi) corrosive fluids (such as brine, oil, and gas) are injected into core samples to measure relative or specific permeability of the core sample. A computer operated permeameter requires a multi-port, three way valve capable of zero volume change and having low internal volume as well as a rapid response time. One prior art approach used VALCO rotary valves such as those manufactured by Valco Instruments, Inc., P.O. Box 55603, Houston, Tex. 77255. These rotary valves, however, have small flow paths and a slow response time such as 0.5-2.0 seconds.\nPrior to the filing of the present invention, a patentability search was conducted resulting in the following patents:\n______________________________________ Inventor Patent No. Date ______________________________________ C. M. Carlson 3,002,532 10-03-61 N. C. Hunt 3,016,917 01-16-62 Maurice F. Franz 3,570,541 03-16-71 Gerald Sorenson 3,587,156 06-28-71 Gerald Sorenson 3,762,443 10-02-73 Stoll et al. 3,794,075 02-26-74 Brake et al. 3,902,526 09-02-75 Charles Perkins 4,027,700 06-07-77 Erich Ruchser 4,067,357 01-10-78 Paul Arvin 4,103,711 08-01-78 Clifford Peters 4,209,040 06-24-80 Clifford Peters 4,217,934 08-19-80 Coppola et al. 4,567,914 02-04-86 ______________________________________\nThe 1971 and 1973 patents to Sorenson pertain to the method of manufacturing a pressurized fluid control valve utilizing a valve spool centrally located in a valve body. The embodiment shown in FIG. 17 of Sorenson shows a multi-passage valve utilizing a symmetrically shaped spool wherein the spool has an incline surface that seats against a sharp circular edge formed in the body. The Sorenson spool or valve body is constructed of deformable resilient materials such as polyurethane, neoprene, polyethylene, or rubber. Deformability of the valve spool or body of Sorenson is critical since in the manufacturing of the valve, the spool or body must deform in order to be inserted into the body. After the insertion, the deformed part recovers to its original condition. In such an embodiment, the forces of deformation are much greater than the forces experienced by the valve in operation. As such, it is believed that the valve of Sorenson would not be suitable to the high temperature, high pressure environments for which the valve of the present invention is ideally suited. Furthermore, Sorenson requires fluid exhaust ports. The present invention has no exhaust ports since it is applied to the flow of high pressure, high temperature fluids which must be completely contained within the valve. Sorenson relies on multiple annular rings providing simultaneous multiple seals caused by moving the spool. The current invention requires only one sealing surface to contact the spool at a time. Sorenson valves are commonly known as 4-way valves, whereas the present valve is specifically a 3-way valve. Finally, the present valve provides means of applying very large forces to a sealing line formed by the poppet (or spool) contacting the sealing surface in the body, namely the large diameter pistons.\nThe 1974 patent to Stoil U.S. Pat. No. (3,794,075) sets forth a multi-way valve having a switching time of six milliseconds and exhibiting short physical valve travel. The Stoll invention makes use of disc plane seals located on a symmetrical center spool. The spool is driven by the snap action of opposing diaphragms which are hydraulically activated. The rubber diaphragms comprising the spool end seals and snap action diaphragms make this valve unsuitable for high pressure applications. High internal pressures would collapse these parts causing the valve to fail. The central O-ring seal of the spool is not constrained so that back pressure on the seal would unseal the O-ring causing leakage. This is not the case in the present invention. Because of flexure of the rubber parts some internal volume change must occur in the Stoll valve after the seals make initial contact. This volume change will compress the fluid in one part of the valve causing a pressure change. This problem does not exist in the present invention. Under high differential pressure the unsupported elastomeric seals used by Stoll may experience extrusion. The problem is accentuated at high temperature. The present valve uses metallic sealing surfaces to eliminate the problem.\nThe 1962 patent issued to Hunt (U.S. Pat. No. 3,016,917) sets forth a valve for controlling high pressure fluid. The valve plunger of Hunt has an enlarged central body portion terminating in an angled seating surface on an enlargement of the valve which selectively abuts against a corresponding angled seat in the cavity. Hunt utilizes a metal to metal engagement between the valve enlargement and the seat to provide a seal against leakage of fluid pressure. The valve is held in the closed position by means of a solenoid and is returned to the opened condition by means of a spring. In this embodiment the central valve spool is not hydrostatically balanced.\nThe 1971 patent to Franz (U.S. Pat. No. 3,570,541) sets forth a three way directional control valve having a formed poppet member located internally to a body wherein the body has internally disposed conically shaped seats for selective engagement with the poppet member to open and close the three way valve. The poppet member and seat incorporate an elastomeric material. The Franz valve is designed to operate at 100 psi with an activating piston driven valve control operating at 70 psi.\nThe 1986 patent to Coppola (U.S. Pat. No. 4,567,914) relates to a two-stage hydraulic solenoid valve having a centrally disposed poppet member having incline surfaces that selectively seat with conical seats. The Coppola invention is designed for steam turbine control systems.\nThe 1978 patent to Ruchser (U.S. Pat. No. 4,067,357) relates to a sliding/spool direction control valve wherein the spool is of symmetrical construction and is driven by two opposing control pistons. The Ruchser embodiment provides a novel seating arrangement that selectively seals despite loss of tolerances due to wear, manufacturing, or use.\nThe remaining patents uncovered in the search are not as pertinent to the present invention as those discussed above."}
{"text": "This invention relates generally to heteroaryl-phenyl substituted compounds and derivatives thereof, which are inhibitors of trypsin-like serine protease enzymes, especially factor Xa, pharmaceutical compositions containing the same, and methods of using the same as anticoagulant agents for treatment and prevention of thromboembolic disorders.\nActivated factor Xa, whose major practical role is the generation of thrombin by the limited proteolysis of prothrombin, holds a central position that links the intrinsic and extrinsic activation mechanisms in the final common pathway of blood coagulation. The generation of thrombin, the final serine protease in the pathway to generate a fibrin clot, from its precursor is amplified by formation of prothrombinase complex (factor Xa, factor V, Ca2+ and phospholipid). Since it is calculated that one molecule of factor Xa can generate 138 molecules of thrombin (Elodi, S., Varadi, K.: Optimization of conditions for the catalytic effect of the factor IXa-factor VIII Complex: Probable role of the complex in the amplification of blood coagulation. Thromb. Res. 1979, 15, 617-629), inhibition of factor Xa may be more efficient than inactivation of thrombin in interrupting the blood coagulation system.\nTherefore, efficacious and specific inhibitors of factor Xa are needed as potentially valuable therapeutic agents for the treatment of thromboembolic disorders. It is thus desirable to discover new factor Xa inhibitors.\nAccordingly, one object of the present invention is to provide novel heteroaryl-phenyl substituted compounds and derivatives thereof that are useful as factor Xa inhibitors or pharmaceutically acceptable salts or prodrugs thereof.\nIt is another object of the present invention to provide pharmaceutical compositions comprising a pharmaceutically acceptable carrier and a therapeutically effective amount of at least one of the compounds of the present invention or a pharmaceutically acceptable salt or prodrug form thereof.\nIt is another object of the present invention to provide a method for treating thromboembolic disorders comprising administering to a host in need of such treatment a therapeutically effective amount of at least one of the compounds of the present invention or a pharmaceutically acceptable salt or prodrug form thereof.\nIt is another object of the present invention to provide novel heteroaryl-phenyl substituted compounds for use in therapy.\nIt is another object of the present invention to provide the use of novel heteroaryl-phenyl substituted compounds for the manufacture of a medicament for the treatment of a thromboembolic disorder.\nThese and other objects, which will become apparent during the following detailed description, have been achieved by the inventors\"\" discovery that the presently claimed heteroaryl-phenyl substituted compounds and derivatives thereof, or pharmaceutically acceptable salt or prodrug forms thereof, are effective factor Xa inhibitors.\n[1] Thus, in a first embodiment, the present invention provides a novel compound of formula Ia, Ib, or Ic: \nor a stereoisomer or pharmaceutically acceptable salt thereof, wherein;\nring D1 is selected from pyridine, pyrazine, pyridazine, and pyrimidine and is substituted with 1 Ra and 0-1 Rb;\nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-3 N atoms, wherein E is selected from O, S, and N-Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nring D3 is a 5-membered heteroaromatic ring system comprising carbon atoms and from 0-3 additional N atoms and ring D3 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, C1-4 alkyl, F, Cl, Br, I, OH, OCH3, OCH2CH3, OCH(CH3)2, OCH2CH2CH3, CN, C(xe2x95x90NR8)NR7R9, NHC(xe2x95x90NR8)NR7R9, NR8CH(xe2x95x90NR7), NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, C(xe2x95x90NH)NH2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), CH2CH2N(C1-3 alkyl)2, (CR8R9)tNR7R8, (CR8R9)tC(O)NR7R8, and OCF3;\nRa is selected from H, C1-4 alkyl, F, Cl, Br, I, OH, OCH3, OCH2CH3, OCH(CH3)2, OCH2CH2CH3, CN, C(xe2x95x90NR8)NR7R9, NHC(xe2x95x90NR8)NR7R9, NR8CH(xe2x95x90NR7), NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, C(xe2x95x90NH)NH2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), CH2CH2N(C1-3 alkyl)2, (CR8R9)tNR7R8, (CR8R9)tC(O)NR7R8, and OCF3;\nRb is selected from H, (C1-4 alkyl, F, Cl, Br, I, OH, OCH3, OCH2CH3, OCH(CH3)2, OCH2CH2CH3, CN, C(xe2x95x90NR8)NR7R9, NHC (xe2x95x90NR8)NR7R9, NR8CH(xe2x95x90NR7), NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, C (xe2x95x90NH)NH2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), CH2CH2N(C1-3 alkyl)2, (CR8R9)tNR7R8, (CR8R9)tC(O)NR7R8, and OCF3;\nRc is selected from H, C1-4 alkyl, OCH3, OCH2CH3, OCH(CH3)2, OCH2CH2CH3, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, C(xe2x95x90NH)NH2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), CH2CH2N(C1-3 alkyl)2, (CR8R9)tNR7R8, (CR8R9)tC(O)NR7R8, and OCF3;\nG is absent or is selected from CH2, C(O), O, NR3, S(O)p, CH2CH2, C(O)CH2, CH2C(O), OCH2, CH2O, NR3CH2, CH2NR3, S(O)pCH2, CH2S(O)p, CH2CH2CH2, C(O)CH2CH2, CH2C(O)CH2, CH2CH2C(O), OCH2CH2, CH2OCH2, CH2CH2O, NR3CH2CH2, CH2NR3CH2, CH2CH2NR3, S(O)pCH2CH2, CH2S(O)pCH2, and CH2CH2S(O)p;\nG1 is absent or is selected from (CR3R3a)1-5, (CR3R3a)0-2CR3xe2x95x90CR3(CR3R3a)0-2(CR3R3a)0-2Cxe2x89xa1C (CR3R3a)0-2, (CR3R3a)uC(O) (CR3R3a)w, (CR3R3a)uC(O)O(CR3R3a)w, (CR3R3a)uOC(O) (CR3R3a)w, (CR3R3a)uO(CR3R3a)w, (CR3R3a)uNR3(CR3R3a)w, (CR3R3a)uC(O)NR3(CR3R3a)w, (CR3R3a)uNR3C(O) (CR3R3a)w, (CR3R3a)uOC(O)NR3(CR3R3a)w, (CR3R3a)uNR3C(O)O(CR3R3a)w, (CR3R3a)uNR3C(O)NR3(CR3R3a)w, (CR3R3a)uNR3C(S)NR3(CR3R3a)w, (CR3R3a)uS (CR3R3a)w, (CR3R3a)uS(O) (CR3R3a)w, (CR3R3a)uS(O)2(CR3R3a)w, (CR3R3a)uS(O)NR3(CR3R3a)w, (CR3R3a)uNR3S(O)2(CR3R3a)w, (CR3R3a)uS(O)2NR3(CR3R3a)w, and (CR3R3a)NR3S(O)2NR3(CR3R3a)w, wherein u+w total 0, 1, 2, 3, or 4, provided that G1 does not form a Nxe2x80x94N, Nxe2x80x94O, Nxe2x80x94S, NCH2N, NCH2O, or NCH2S bond with either group to which it is attached;\nG2 is phenyl, naphthyl, or a 5-10 membered heteroaryl consisting of carbon atoms and from 1-3 heteroatoms selected from N, O, and S;\nM is isoxazoline, pyrazoline, isothiazoline, triazoline, tetrazoline, phenyl, or a 5-6 membered aromatic heterocycle consisting of carbon atoms and 1-4 heteroatoms selected from O, N, and S, and is substituted with -Z-A-B;\nM is also substituted with 0-2 R1a;\nZ is selected from a bond, xe2x80x94(CR2R2a)1-4xe2x80x94, (CR2R2a)qO(CR2R2a)q1, (CR2R2a)qNR3(CR2R2a)q1, (CR2R2a)qC(O) (CR2R2a)q1, (CR2R2a)qC(O)O(CR2R2a)q1, (CR2R2a)qOC(O) (CR2R2a)q1(CR2R2a)qC(O)NR3(CR2R2a)q1, (CR2R2a)qNR3C(O) (CR2R2a)q1, (CR2R2a)qOC(O)O(CR2R2a)q1, (CR2R2a)qOC(O)NR3 (CR2R2a)q1, (CR2R2a)qNR3C(O)O(CR2R2a)q1, (CR2R2a)q1NR3C(O) NR3(CR2a)q1, (CR2R2a)qS(CR2R2a)q1, (CR2R2a)qS(O) (CR2R2a)q1, (CR2R2a)qS(O)2(CR2R2a)q1, (CR2R2a)qSO2NR3(CR2R2a)q1, (CR2R2a)qNR3SO2(CR2R2a)q1, and (CR2R2a)qNR3SO2NR3(CR2R2a)q1, wherein q+q1 total 0, 1, or 2, provided that Z does not form a Nxe2x80x94N, Nxe2x80x94O, Nxe2x80x94S, NCH2N, NCH2O, or NCH2S bond with either group to which it is attached;\nR1a is selected from H, xe2x80x94(CH2)rxe2x80x94R1b, xe2x80x94CH=CHxe2x80x94R1b, NCH2R1c, OCH2R1c, SCH2R1c, NH(CH2)2(CH2)tR1b, O(CH2)2(CH2)tR1b, S(CH2)2(CH2)tR1b, S(O)p(CH2)rR1d, O(CH2)rR1d, NR3 (CH2)rR1d, OC(O)NR3(CH2)rR1d, NR3C(O)NR3(CH2)rR1d, NR3C(O)O(CH2)rR1d, and NR3C(O) (CH2)rR1d, provided that R1a forms other than an Nxe2x80x94halo, Nxe2x80x94N, Nxe2x80x94S, Nxe2x80x94O, or Nxe2x80x94CN bond;\nalternatively, when two R1a\"\"s are attached to adjacent atoms, together with the atoms to which they are attached they form a 5-7 membered ring consisting of: carbon atoms and 0-2 heteroatoms selected from the group consisting of N, O, and S(O)p, this ring being substituted with 0-2 R4b and comprising: 0-3 double bonds;\nR1b is selected from H, C1-3 alkyl, F, Cl, Br, I, xe2x80x94CN, xe2x80x94CHO, (CF2)rCF3, (CH2)rOR2, NR2R2a, C(O)R2c, OC(O)R2, (CF2)rCO2R2a, S(O)pR2b, NR2(CH2)rOR2, C(xe2x95x90NR2c)NR2R2a, NR2C(O)R2b, NR2C(O)NHR2b, NR2C(O)2R2a, OC(O)NR2aR2b, C(O)NR2R2a, C(O)NR2(CH2)rOR2, SO2NR2R2a, NR2SO2R2b, C3-6 carbocycle substituted with 0-2 R4a, and 5-10 membered heterocycle consisting of carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-2 R4a, provided that R1b forms other than an Nxe2x80x94halo, Nxe2x80x94N, Nxe2x80x94S, Nxe2x80x94O, or Nxe2x80x94CN bond;\nR1c is selected from H, CH(CH2OR2)2, C(O)R2c, C(O)NR2R2a, S(O)R2b, S(O)2R2b, and SO2NR2R2a;\nR1d is selected from C3-13 carbocycle substituted with 0-2 R4a, and 5-13 membered heterocycle consisting of carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-2 R4a, provided that R1d forms other than an Nxe2x80x94N, Nxe2x80x94S, or Nxe2x80x94O bond;\nR2, at each occurrence, is selected from H, CF3, C1-6 alkyl, benzyl, C3-6 carbocyclic residue substituted with 0-2 R4b, and 5-6 membered heterocyclic system comprising carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S substituted with 0-2 R4b;\nR2a, at each occurrence, is selected from H, CF3, C1-6 alkyl, benzyl, phenethyl, C3-6 carbocyclic residue substituted with 0-2 R4b, and 5-6 membered heterocyclic system comprising carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S substituted with 0-2 R4b;\nR2b, at each occurrence, is selected from CF3, C1-4 alkoxy, C1-6 alkyl, benzyl, C3-6 carbocyclic residue substituted with 0-2 R4b, and 5-6 membered heterocyclic system comprising carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S substituted with 0-2 R4b;\nR2c, at each occurrence, is selected from CF3, OH, C1-4 alkoxy, C1-6 alkyl, benzyl, C3-6 carbocyclic residue substituted with 0-2 R4b, and 5-6 membered heterocyclic system comprising carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S substituted with 0-2 R4b;\nalternatively, R2 and R2a, together with the atom to which they are attached, combine to form a 5 or 6 membered saturated, partially saturated or unsaturated ring substituted with 0-2 R4b and comprising carbon atoms and from 0-1 additional heteroatoms selected from the group consisting of N, O, and S;\nR3, at each occurrence, is selected from H, C1-4 alkyl, and phenyl;\nR3a, at each occurrence, is selected from H, C1-4 alkyl, and phenyl;\nR3b, at each occurrence, is selected from H, C1-4 alkyl, and phenyl;\nR3c, at each occurrence, is selected from C1-4 alkyl, and phenyl;\nR3d, at each occurrence, is selected from H, (C1-4 alkyl, C1-4 alkyl-phenyl, and C(xe2x95x90O)R3c;\nA is selected from:\nC3-10 carbocyclic residue substituted with 0-2 R4, and\n5-12 membered heterocyclic system comprising carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S substituted with 0-2 R4;\nB is selected from: H, Y, and X-Y, provided that Z and B are attached to different atoms on A;\nX is selected from xe2x80x94(CR2R2a)14xe2x80x94, xe2x80x94R2(CR2R2b) (CH2)txe2x80x94, xe2x80x94C(O)xe2x80x94, xe2x80x94C(xe2x95x90NR1c)xe2x80x94, xe2x80x94CR2(NR1cR2)xe2x80x94, xe2x80x94CR2(OR2)xe2x80x94, xe2x80x94CR2(SR2)xe2x80x94, xe2x80x94C(O)CR2R2axe2x80x94, xe2x80x94CR2R2aC(O), xe2x80x94Sxe2x80x94, xe2x80x94S(O)xe2x80x94, xe2x80x94S(O)2xe2x80x94, xe2x80x94SCR2R2axe2x80x94, xe2x80x94S(O)CR2R2axe2x80x94, xe2x80x94S(O)2CR2R2axe2x80x94, xe2x80x94CR2R2aSxe2x80x94, xe2x80x94CR2R2aS(O)xe2x80x94, xe2x80x94CR2R2aS(O)2xe2x80x94, xe2x80x94S(O)2NR2xe2x80x94, xe2x80x94NR2S(O)2xe2x80x94, xe2x80x94NR2S(O)2CR2R2axe2x80x94, xe2x80x94CR2R2aS(O)2NR2xe2x80x94, xe2x80x94NR2S(O) 2NR2xe2x80x94, xe2x80x94C(O)NR2xe2x80x94, xe2x80x94NR2C(O)xe2x80x94, xe2x80x94C(O)NR2CR2R2axe2x80x94, xe2x80x94NR2C(O)CR2R2axe2x80x94, xe2x80x94CR2R2aC(O)NR2xe2x80x94, xe2x80x94CR2R2aNR2C(O)xe2x80x94, xe2x80x94NR2C(O)Oxe2x80x94, xe2x80x94OC(O)NR2xe2x80x94, xe2x80x94NR2C(O)NR2xe2x80x94, xe2x80x94NR2xe2x80x94, xe2x80x94NR2CR2R2axe2x80x94, xe2x80x94CR2R2aNR2xe2x80x94, O, xe2x80x94CR2R2aOxe2x80x94, and xe2x80x94OCR2R2axe2x80x94;\nY is selected from:\nC3-10 carbocyclic residue substituted with 0-2 R4a, and\n5-12 membered heterocyclic system comprising carbon atoms and from 1-4 heteroatoms selected from the group consisting of N, O, and S substituted with 0-2 R4a;\nR4, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR2, (CH2)rF, (CH2)rCl, (CH2)rBr, (CH2)rI, C1-4 alkyl, (CH2)rCN, (CH2)rNO2, (CH2)rNR2R2a, C(O)R2c, NR2C(O)R2b, C(O)NR2R2a, NR2C(O)NR2R2a, C(xe2x95x90NR2)NR2R2a, C(xe2x95x90NS(O)2R5)NR2R2a, NHC(xe2x95x90NR2)NR2R2a, C(O)NHC(xe2x95x90NR2)NR2R2a, SO2NR2R2a, NR2SO2NR2R2a, NR2SO2xe2x80x94C1-4 alkyl, NR2SO2R5, S(O)pR5, (CF2)rCF3, (CH2)rxe2x80x94CF3, NCH2R1c, OCH2R1c, SCH2R1c, N(CH2)2(CH2)tR1b, O(CH2)2(CH2)tR1b, S(CH2)2(CH2)tR1b, 5-6 membered carbocycle substituted with 0-1 R5, and 5-6 membered heterocycle consisting of: carbon atoms and 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-1 R5;\nR4a, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR2, (CF2)rCF3, (CH2)rxe2x80x94CF3, (CH2)rxe2x80x94F, (CH2)rxe2x80x94Br, (CH2)rxe2x80x94Cl, (C1-4 alkyl, (CH2)rCN, (CH2)rNO2, (CH2)rNR2R2a, (CH2)rC(O)R2c, NR2C(O)R2b, C(O)NR2R2a, (CH2)rNxe2x80x94CHOR3, C(O)NH(CH2)2NR2R2a, NR2C(O)NR2R2a, C(xe2x95x90NR2)NR2R2a, NHC(xe2x95x90NR2)NR2R2a, SO2NR2R2a, NR2SO2NR2R2a, NR2SO2xe2x80x94C1-4 alkyl, NR2SO2R5, C(O)NHSO2xe2x80x94C1-4 alkyl, S(O)pR5, 5-6 membered carbocycle substituted with 0-1 R5, and 5-6 membered heterocycle consisting of: carbon atoms and 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-1 R5;\nR4b, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR3, (CH2)rxe2x80x94F, (CH2)rxe2x80x94Cl, (CH2)rxe2x80x94Br, (CH2)rxe2x80x94I, (C1-4 alkyl, (CH2)rxe2x80x94CN, (CH2)rxe2x80x94NO2, (CH2)rNR3R3a, (CH2)rC(O)R3, (CH2)rC(O)OR3c, NR3C(O)R3a, C(O)NR3R3a, NR3C(O)NR3R3a, C(xe2x95x90NR3)NR3R3a, NR3C(xe2x95x90NR3)NR3R3a, SO2NR3R3a, NR3SO2NR3R3a, NR3SO2xe2x80x94C1-4 alkyl, NR3SO2CF3, NR3SO2-phenyl, S(O)pCF3, S(O)pxe2x80x94C1-4 alkyl, S(O)p-phenyl, (CH2)rCF3, and (CF2)rCF3;\nR5, at each occurrence, is selected from H, C1-6 alkyl, xe2x95x90O, (CH2)rOR3, F, Cl, Br, I, xe2x80x94CN, NO2, (CH2)rNR3R3a, (CH2)rC(O)R3, (CH2)rC(O)OR3c, NR3C(O)R3a, C(O)NR3R3a, NR3C(O)NR3R3a, CH(xe2x95x90NOR3d), C(xe2x95x90NR3)NR3R3a, NR3C(xe2x95x90NR3)NR3R3a, SO2NR3R3a, NR3SO2NR3R3a, NR3SO2xe2x80x94C1-4 alkyl, NR3SO2CF3, NR3SO2-phenyl, S(O)pCF3, S(O)pxe2x80x94C1-4 alkyl, S(O)p-phenyl, (CF2)rCF3, phenyl substituted with 0-2 R6, naphthyl substituted with 0-2 R6, and benzyl substituted with 0-2 R6;\nR6, at each occurrence, is selected from H, OH, (CH2)rOR2, halo, C1-4 alkyl, CN, NO2, (CH2)rNR2R2a, (CH2)rC(O)R2b, NR2C(O)R2b, NR2C(O)NR2R2a, C(xe2x95x90NH)NH2, NHC(xe2x95x90NH)NH2, SO2NR2R2a, NR2SO2NR2R2a, and NR2SO2C1-4 alkyl;\nR7, at each occurrence, is selected from H, OH, C1-4 alkoxycarbonyl, C6-10 aryloxy, C6-10 aryloxycarbonyl, C6-10 arylmethylcarbonyl, C1-4 alkylcarbonyloxy C1-4 alkoxycarbonyl, C6-10 arylcarbonyloxy C1-4 alkoxycarbonyl, C1-6 alkylaminocarbonyl, phenylaminocarbonyl, and phenyl (C1-4 alkoxycarbonyl;\nR8, at each occurrence, is selected from H, C1-6 alkyl, and (CH2)n-phenyl;\nalternatively, R7 and R8, when attached to the same nitrogen, combine to form a 5-6 membered heterocyclic ring consisting of carbon atoms and 0-2 additional heteroatoms selected from the group consisting of N, O, and S(O)p;\nR9, at each occurrence, is selected from H, C1-6 alkyl, and (CH2)n-phenyl;\nn, at each occurrence, is selected from 0, 1, 2, and 3;\nm, at each occurrence, is selected from 0, 1, and 2;\np, at each occurrence, is selected from 0, 1, and 2;\nr, at each occurrence, is selected from 0, 1, 2, and 3;\ns, at each occurrence, is selected from 0, 1, and 2; and,\nt, at each occurrence, is selected from 0, 1, 2, and 3.\n[2] In a preferred embodiment, the present invention provides a novel compound, wherein the compound is of formula Ia1-Ic1, wherein: \nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-2 N atoms, wherein E is selected from O, S, and Nxe2x80x94Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nring D3 is a 5-membered heteroaromatic ring system comprising carbon atoms and from 0-3 additional N atoms and ring D3 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N(C1-3 alkyl)2;\nRa is selected from H, OH, SH, C1-3 alkoxy, C1-3 thioalkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N (C1-3 alkyl)2,\nRb is selected from H, C1-4 alkyl, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N(C1-3 alkyl)2;\nRc is selected from H, C1-4 alkyl, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N(C1-3 alkyl)2;\nG1 is absent or is selected from CH2, C(O), O, NR3, S(O)p, CH2CH2, C(O)CH2, CH2C(O), OCH2, CH2O, NR3CH2, CH2NR3, S(O)pCH2, CH2S(O)p, CH2CH2CH2, C(O)CH2CH2, CH2C(O)CH2, CH2CH2C(O), OCH2CH2, CH2OCH2, CH2CH2O, NR3CH2CH2, CH2NR3CH2, CH2CH2NR3, S(O)pCH2CH2, CH2S(O)pCH2, and CH2CH2S(O)p, and provided that G1xe2x80x94M form other than a Nxe2x80x94N, Oxe2x80x94N, or Sxe2x80x94N bond;\nM is selected from the group: \nJ is O or S;\nJa is NH or NR1a;\nG1 is absent or is selected from (CR3R3a)1-3, (CR3R3a)uC(O) (CR3R3a)w, (CR3R3a)uO(CR3R3a)w, (CR3R3a)uNR3(CR3R3a)w, (CR3R3a)uC(O)NR3(CR3R3a)w, (CR3R3a)uNR3C(O) (CR3R3a)w, (CR3R3a)uS(CR3R3a)w, (CR3R3a)uS(O) (CR3R3a)w, (CR3R3a) S(O)2(CR3R3a)w, (CR3R3a)uS(O)NR3CR3R3a)w, and (CR3R3a)uS(O)2NR3(CR3R3a)w, wherein u+w total 0, 1, or 2, provided that G1 does not form a Nxe2x80x94N, Nxe2x80x94O, Nxe2x80x94S, NCH2N, NCH2O, or NCH2S bond with either group to which it is attached;\nA is selected from one of the following carbocyclic and heterocyclic systems which are substituted with 0-2 R4;\nphenyl, piperidinyl, piperazinyl, pyridyl, pyrimidyl, furanyl, morpholinyl, thiophenyl, pyrrolyl, pyrrolidinyl, oxazolyl, isoxazolyl, thiazolyl, isothiazolyl, pyrazolyl, imidazolyl, oxadiazolyl, thiadiazolyl, triazolyl, 1,2,3-oxadiazolyl, 1,2,4-oxadiazolyl, 1,2,5-oxadiazolyl, 1,3,4-oxadiazolyl, 1,2,3-thiadiazolyl, 1,2,4-thiadiazolyl, 1,2,5-thiadiazolyl, 1,3,4-thiadiazolyl, 1,2,3-triazolyl, 1,2,4-triazolyl, 1,2,5-triazolyl, 1,3,4-triazolyl, benzofuranyl, benzothiofuranyl, indolyl, benzimidazolyl, benzoxazolyl, benzthiazolyl, indazolyl, benzisoxazolyl, benzisothiazolyl, and isoindazolyl;\nX is selected from xe2x80x94C(CR2R2a)14xe2x80x94, xe2x80x94C(O)xe2x80x94, xe2x80x94C(xe2x95x90NR1c)xe2x80x94, xe2x80x94CR2(NR1cR2)xe2x80x94C(O)CR2R2axe2x80x94, xe2x80x94CR2R2aC(O)xe2x80x94C(O)NR2xe2x80x94, xe2x80x94NR2C(O)xe2x80x94, C(O)NR2CR2R2axe2x80x94, xe2x80x94NR2C(O)CR2R2axe2x80x94, xe2x80x94CR2R2aC(O)NR2xe2x80x94, xe2x80x94CR2R2aNR2C(O)xe2x80x94, xe2x80x94NR2C(O)NR2xe2x80x94, xe2x80x94NR2xe2x80x94, xe2x80x94NR2CR2R2axe2x80x94, xe2x80x94CR2R2aNR2xe2x80x94, O, xe2x80x94CR2R2aOxe2x80x94, and xe2x80x94OCR2R2axe2x80x94;\nY is selected from one of the following carbocyclic and heterocyclic systems that are substituted with 0-2 R4a;\ncyclopropyl, cyclopentyl, cyclohexyl, phenyl, piperidinyl, piperazinyl, pyridyl, pyrimidyl, furanyl, morpholinyl, thiophenyl, pyrrolyl, pyrrolidinyl, oxazolyl, isoxazolyl, isoxazolinyl, thiazolyl, isothiazolyl, pyrazolyl, imidazolyl, oxadiazolyl, thiadiazolyl, triazolyl, 1,2,3-oxadiazolyl, 1,2,4-oxadiazolyl, 1,2,5-oxadiazolyl, 1,3,4-oxadiazolyl, 1,2,3-thiadiazolyl, 1,2,4-thiadiazolyl, 1,2,5-thiadiazolyl, 1,3,4-thiadiazolyl, 1,2,3-triazolyl, 1,2,4-triazolyl, 1,2,5-triazolyl, 1,3,4-triazolyl, benzofuranyl, benzothiofuranyl, indolyl, benzimidazolyl, benzoxazolyl, benzthiazolyl, indazolyl, benzisoxazolyl, benzisothiazolyl, and isoindazolyl;\nalternatively, Y is selected from the following bicyclic heteroaryl ring systems: \nK is selected from O, S, NH, and N;\nZ is selected from a bond, CH2O, OCH2, NH, CH2NH, NHCH2, CH2C(O)C(O)CH2, C(O) NH, NHC(O), CH2S(O)2, S(O)2(CH2), SO2NH, and NHSO2, provided that Z does not form a Nxe2x80x94N, Nxe2x80x94O, Nxe2x80x94S, NCH2N, NCH2O, or NCH2S bond with either group to which it is attached;\nR4, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR2, F, Cl, Br, I, C1-4 alkyl, CN, NO2, (CH2)rNR2R2a, C(O)R2c, NR2C(O)R2b, C(O)NR2R2a, NR2C(O)NR2R2a, C(xe2x95x90NR2)NR2R2a, SO2NR2R2a, NR2SO2NR2R2a, NR2SO2xe2x80x94C1-4 alkyl, NR2SO2R5, S(O)pR5, CF3, NCH2R1c, OCH2R1c, SCH2R1c, N(CH2)2(CH2)tR1b, O(CH2)2(CH2)tR1b, S(CH2)2(CH2)tR1b, 5-6 membered carbocycle substituted with 0-1 R5, and 5-6 membered heterocycle consisting of: carbon atoms and 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-1 R5; and,\nR4a, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR2, CF3, F, Br, Cl, C1-4 alkyl, CN, NO2, (CH2)rNR2R2a, (CH2)rC(O)R2c, NR2C(O)R2b, C(O)NR2R2a, NR2C(O)NR2R2a, C(xe2x95x90NR2)NR2R2a, NHC(xe2x95x90NR2)NR2R2a, SO2NR2R2a, NR2SO2NR2R2a, NR2SO2xe2x80x94C1-4 alkyl, NR2SO2R5, C(O)NHSO2xe2x80x94C1-4 alkyl, S(O)pR5, 5-6 membered carbocycle substituted with 0-1 R5, and 5-6 membered heterocycle consisting of: carbon atoms and 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-1 R5.\n[3] In another preferred embodiment, the present invention provides a novel compound, wherein the compound is of formula Ib1 or Ic1, wherein;\nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-2 N atoms, wherein E is selected from O, S, and Nxe2x80x94Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nRa is selected from H. OH, SH, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nRb is selected from H, C1-4 alkyl, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nRc is selected from H, C1-4 alkyl, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N (C1-3 alkyl)2;\nM is selected from the group: \nY is selected from one of the following carbocyclic and heterocyclic systems which are substituted with 0-2 R4a;\nphenyl, piperidinyl, piperazinyl, pyridyl, pyrimidyl, furanyl, morpholinyl, thiophenyl, pyrrolyl, pyrrolidinyl, oxazolyl, isoxazolyl, thiazolyl, isothiazolyl, pyrazolyl, imidazolyl, oxadiazole, thiadiazole, triazole, 1,2,3-oxadiazole, 1,2,4-oxadiazole, 1,2,5-oxadiazole, 1,3,4-oxadiazole, 1,2,3-thiadiazole, 1,2,4-thiadiazole, 1,2,5-thiadiazole, 1,3,4-thiadiazole, 1,2,3-triazole, 1,2,4-triazole, 1,2,5-triazole, 1,3,4-triazole, benzofuran, benzothiofuran, indole, benzimidazole, benzimidazolone, benzoxazole, benzthiazole, indazole, benzisoxazole, benzisothiazole, and isoindazole;\nZ is selected from a bond, CH2O, OCH2, NH, CH2NH, NHCH2, CH2C(O), C(O)CH2, C(O)NH, NHC(O), CH2S(O)2, S(O)2(CH2), SO2NH, and NHSO2, provided that Z does not form a Nxe2x80x94N, Nxe2x80x94O, Nxe2x80x94S, NCH2N, NCH2O, or NCH2S bond with either group to which it is attached;\nR4, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR2, F, Cl, Br, I, C1-4 alkyl, CN, NO2, (CH2)rNR2R2a, C(O)R2c, NR2C(O)R2b, C(O)NR2R2a, NR2C(O)NR2R2a, C(xe2x95x90NR2)NR2R2a, SO2NR2R2a, NR2SO2NR2R2a, NR2SO2xe2x80x94C1-4 alkyl, NR2SO2R5, S(O)pR5, CF3, 5-6 membered carbocycle substituted with 0-1 R5, and 5-6 membered heterocycle consisting of: carbon atoms and 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-1 R5; and,\nR4a, at each occurrence, is selected from H, xe2x95x90O, (CH2)rOR2, CF3, F, Br, Cl, C1-4 alkyl, CN, NO2, (CH2)rNR2R2a, (CH2)rC(O)R2c, NR2C(O)R2b, C(O)NR2R2a, NR2C(O)NR2R2a, C(xe2x95x90NR2)NR2R2a, SO2NR2R2a, C(O)NHSO2xe2x80x94C1-4 alkyl, S(O)pR5, 5-6 membered carbocycle substituted with 0-1 R5, and 5-6 membered heterocycle consisting of: carbon atoms and 1-4 heteroatoms selected from the group consisting of N, O, and S(O)p substituted with 0-1 R5.\n[4] In another preferred embodiment, the present invention provides a novel compound, wherein the compound is of formula Ib2: \nor a stereoisomer or pharmaceutically acceptable salt thereof, wherein;\nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-2 N atoms, wherein E is selected from O, S, and Nxe2x80x94Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nRa is selected from H, OH, SH, NH2, NH(C1-3 alkyl), and N(C1-3 alkyl)2;\nRb is selected from H, C1-4 alkyl, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nM is selected from the group: \nG1 is absent or is selected from CH2, CH2CH2, CH2O, OCH2, NH, CH2NH, NHCH2 CH2C(O), C(O)CH2, C(O)NH, NHC(O), CH2S(O)2, S(O)2(CH2), SO2NH, and NHSO2, provided that G1 does not form a Nxe2x80x94N, Nxe2x80x94O, Nxe2x80x94S, NCH2N, NCH2O, or NCH2S bond with either group to which it is attached.\n[5] In another preferred embodiment, the present invention provides a novel compound, wherein the compound is selected from one of the formulas: \nor a stereoisomer or pharmaceutically acceptable salt thereof, wherein;\nG1 is absent;\nA is selected from phenyl, piperidinyl, pyridyl, and\npyrimidyl, and is substituted with 0-2 R4; and,\nB is selected from phenyl, pyrrolidino, N-pyrrolidino-carbonyl, morpholino, N-morpholino-carbonyl, 1,2,3-triazolyl, imidazolyl, and benzimidazolyl, and is substituted with 0-1 R4a;\nR2, at each occurrence, is selected from H, CH3, CH2CH3, cyclopropylmethyl, cyclobutyl, and cyclopentyl;\nR2a, at each occurrence, is selected from H, CH3, and CH2CH3;\nalternatively, R2 and R2a, together with the atom to which they are attached, combine to form pyrrolidine substituted with 0-2 R4b or piperidine substituted with 0-2 R4b;\nR4, at each occurrence, is selected from OH, OR2, (CH2)OR2, (CH2)2OR2, F, Br, Cl, I, C1-4 alkyl, NR2R2a, (CH2)NR2R2a, (CH2)2NR2R2a, CF3, and (CF2)CF3;\nR4a is selected from C1-4 alkyl, CF3, OR2, (CH2)OR2, (CH2)2OR2, NR2R2a, (CH2)NR2R2a, (CH2)2NR2R2a, SR5, S(O)R5, S(O)2R5, SO2NR2R2a, and 1xe2x80x94CF3-tetrazol-2-yl;\nR4b, at each occurrence, is selected from H, CH3, and OH;\nR5, at each occurrence, is selected from CF3, C1-6 alkyl, phenyl, and benzyl; and,\nr, at each occurrence, is selected from 0, 1, and 2.\n[6] In another preferred embodiment, the present invention provides a novel compound, wherein;\nA is selected from the group: phenyl, piperidinyl, 2-pyridyl, 3-pyridyl, 2-pyrimidyl, 2xe2x80x94Cl-phenyl, 3-Cl-phenyl, 2-F-phenyl, 3-F-phenyl, 2-methylphenyl, 2-aminophenyl, and 2-methoxyphenyl; and,\nB is selected from the group: 2-(aminosulfonyl)phenyl, 2-(methylaminosulfonyl)phenyl, 1-pyrrolidinocarbonyl, 2-(methylsulfonyl)phenyl, 2-(N,N-dimethylaminomethyl)phenyl, 2-(N-methylaminomethyl)phenyl, 2-(N-ethyl-N-methylaminomethyl)phenyl, 2-(N-pyrrolidinylmethyl)phenyl, 1-methyl-2-imidazolyl, 2-methyl-1-imidazolyl, 2-(dimethylaminomethyl)-1-imidazolyl, 2-(methylaminomethyl)-1-imidazolyl, 2-(N-(cyclopropylmethyl)aminomethyl)phenyl, 2-(N-(cyclobutyl)aminomethyl)phenyl, 2-(N-(cyclopentyl)aminomethyl)phenyl, 2-(N-(4-hydroxypiperidinyl)methyl)phenyl, and 2-(N-(3-hydroxypyrrolidinyl)methyl)phenyl.\n[7] In another preferred embodiment, the present invention provides a novel compound selected from the group:\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(2xe2x80x2-N,N-dimethylaminomethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(2xe2x80x2-N-pyrrolidinylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N,N-dimethylaminomethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N-pyrrolidinylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(2xe2x80x2-sulfonylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-sulfonylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(2xe2x80x2-Amino-1xe2x80x2,3xe2x80x2,4xe2x80x2-triazol-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-sulfonylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(2xe2x80x2-Amino-1xe2x80x2,3xe2x80x2,4xe2x80x2-triazol-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N,N-dimethylaminomethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(2xe2x80x2-Amino-1xe2x80x2,3xe2x80x2,4xe2x80x2-triazol-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N-pyrrolidinylmethyl-[1,1xe2x80x2],-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(2xe2x80x2-Amino-3xe2x80x2,4xe2x80x2-thiadiazol-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-sulfonylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(2xe2x80x2-Amino-3xe2x80x2,4xe2x80x2-thiadiazol-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N,N-dimethylaminomethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(2xe2x80x2-Amino-3xe2x80x2,4xe2x80x2-thiadiazol-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N-pyrrolidinylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[4-(2xe2x80x2-(N,N-dimethylamino)methylimidazol-1xe2x80x2-yl)-2-fluorophenyl)aminocarbonyl]pyrazole;\n1-[3-(1xe2x80x2,3xe2x80x2,4xe2x80x2-Triazol-2xe2x80x2-on-5xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[4-(2xe2x80x2-methylimidazol-1xe2x80x2-yl)-2-fluorophenyl)aminocarbonyl]pyrazole;\n1-[3-(Pyrid-3xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-sulfonylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(Pyrid-3xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N,N-dimethylaminomethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(Pyrid-3xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N-pyrrolidinylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(Pyrid-4xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N,N-dimethylaminomethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(Pyrid-4xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-N-pyrrolidinylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(Pyrid-4xe2x80x2-yl)phenyl]-3-trifluoromethyl-5-[(3-fluoro-2xe2x80x2-sulfonylmethyl-[1,1xe2x80x2]-biphen-4-yl)aminocarbonyl]pyrazole;\n1-[3-(5-Amino-1,3,4-oxadiazol-2-yl)phenyl]-N-{2xe2x80x2-[(dimethylamino)methyl]-3-fluoro-1,1xe2x80x2-biphenyl-4-yl}-3-(trifluoromethyl)-1H-pyrazole-5-carboxamide;\n1-[3-(5-Amino-1,3,4-oxadiazol-2-yl)phenyl]-N-(3-fluoro-2xe2x80x2-{[(3S)-3-hydroxy-1-pyrrolidinyl]methyl}-1,1xe2x80x2-biphenyl-4-yl)-3-(trifluoromethyl)-1H-pyrazole-5-carboxamide; and,\n1-[3-(5-Amino-1,3,4-oxadiazol-2-yl)phenyl]-N-(3-fluoro-2xe2x80x2-{[(3S)-3-hydroxy-1-pyrrolidinyl]methyl}-1,1xe2x80x2-biphenyl-4-yl)-5-(trifluoromethyl)-1H-pyrazole-3-carboxamide;\nor a pharmaceutically acceptable salt thereof.\nIn another embodiment, the present invention provides a novel compound, wherein the compound is of formula Ia1-Ic1, wherein: \nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-2 N atoms, wherein E is selected from O, S, and Nxe2x80x94Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nring D3 is a 5-membered heteroaromatic ring system comprising carbon atoms and from 0-3 additional N atoms and ring D3 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N(C1-3 alkyl)2;\nRa is selected from H, OH, SH, C1-3 alkoxy, C1-3 thioalkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N(C1-3 alkyl)2;\nRb is selected from H, C1-4 alkyl, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N(C1-3 alkyl)2;\nRc is selected from H, C1-4 alkyl, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), CH2N(C1-3 alkyl)2, CH2CH2NH2, CH2CH2NH(C1-3 alkyl), and CH2CH2N (C1-3 alkyl)2, and,\nG1 is absent or is selected from CH2, C(O), O, NR3, S(O)p, CH2CH2, C(O)CH2, CH2C(O), OCH2, CH2O, NR3CH2, CH2NR3, S(O)pCH2, CH2S(O)p, CH2CH2CH2, C(O)CH2CH2, CH2C(O)CH2, CH2CH2C(O), OCH2CH2, CH2OCH2, CH2CH2O, NR3CH2CH2, CH2NR3CH2, CH2CH2NR3, S(O)pCH2CH2, CH2S(O)CH2, and CH2CH2S(O)p, and provided that G1xe2x80x94M form other than a Nxe2x80x94N, Oxe2x80x94N, or Sxe2x80x94N bond.\nIn another embodiment, the present invention provides a novel compound, wherein the compound is of formula Ib1 or Ic1, wherein;\nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-2 N atoms, wherein E is selected from O, S, and Nxe2x80x94Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nRa is selected from H, OH, SH, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2,\nRb is selected from H, C1-4 alkyl, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2; and,\nRc is selected from H, C1-4 alkyl, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2.\nIn another embodiment, the present invention provides a novel compound, wherein the compound is of formula Ib2: \nor a stereoisomer or pharmaceutically acceptable salt thereof, wherein;\nring D2 is a 5-membered heteroaromatic ring system comprising E, carbon atoms, and 0-2 N atoms, wherein E is selected from O, S, and Nxe2x80x94Rc and ring D2 is substituted with 1 Ra and 0-1 Rb;\nR is selected from H, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2;\nRa is selected from H. OH, SH, NH2, NH(C1-3 alkyl), and N(C1-3 alkyl)2; and,\nRb is selected from H, C1-4 alkyl, Cl, F, Br, I, OH, C1-3 alkoxy, NH2, NH(C1-3 alkyl), N(C1-3 alkyl)2, CH2NH2, CH2NH(C1-3 alkyl), and CH2N(C1-3 alkyl)2.\nIn another embodiment, the present invention provides a novel compound selected from one of the formulas: \nor a stereoisomer or pharmaceutically acceptable salt thereof.\nIn another embodiment, the present invention provides a novel compound selected from one of the formulas: \nor a stereoisomer or pharmaceutically acceptable salt thereof.\nIn another embodiment, the present invention provides a novel compound selected from one of the formulas: \nor a stereoisomer or pharmaceutically acceptable salt thereof.\nIn another embodiment, the present invention provides a novel compound wherein M is selected from the group: \nJ is O or S; and,\nJa is NH or NR1a.\nIn another embodiment, the present invention provides a novel compound wherein M is selected from the group: \nIn another embodiment, the present invention provides a novel compound wherein M is selected from the group: \nIn another embodiment, the present invention provides a novel compound wherein A is selected from one of the following carbocyclic and heterocyclic systems that are substituted with 0-2 R4;\nphenyl, piperidinyl, piperazinyl, pyridyl, pyrimidyl, furanyl, morpholinyl, thiophenyl, pyrrolyl, pyrrolidinyl, oxazolyl, isoxazolyl, thiazolyl, isothiazolyl, pyrazolyl, imidazolyl, oxadiazolyl, thiadiazolyl, triazolyl, 1,2,3-oxadiazolyl, 1,2,4-oxadiazolyl, 1,2,5-oxadiazolyl, 1,3,4-oxadiazolyl, 1,2,3-thiadiazolyl, 1,2,4-thiadiazolyl, 1,2,5-thiadiazolyl, 1,3,4-thiadiazolyl, 1,2,3-triazolyl, 1,2,4-triazolyl, 1,2,5-triazolyl, 1,3,4-triazolyl, benzofuranyl, benzothiofuranyl, indolyl, benzimidazolyl, benzoxazolyl, benzthiazolyl, indazolyl, benzisoxazolyl, benzisothiazolyl, and isoindazolyl;\nIn another embodiment, the present invention provides a novel compound wherein A is selected from phenyl, piperidinyl, pyridyl, and pyrimidyl, and is substituted with 0-2 R4.\nIn another embodiment, the present invention provides a novel compound wherein A is selected from the group: phenyl, piperidinyl, 2-pyridyl, 3-pyridyl, 2-pyrimidyl, 2-Cl-phenyl, 3-Cl-phenyl, 2-F-phenyl, 3-F-phenyl, 2-methylphenyl, 2-aminophenyl, and 2-methoxyphenyl.\nIn another embodiment, the present invention provides a novel compound wherein:\nB is selected from: H, Y, and X-Y, provided that Z and B are attached to different atoms on A;\nX is selected from xe2x80x94(CR2R2a)1-4, xe2x80x94C(O)xe2x80x94C(xe2x95x90NR1c)xe2x80x94, xe2x80x94CR2(NR1cR2), xe2x80x94C(O)CR2R2axe2x80x94, xe2x80x94CR2R2aC(O), xe2x80x94C(O)NR2xe2x80x94, xe2x80x94NR2C(O)xe2x80x94, xe2x80x94C(O)NR2CR2R2axe2x80x94, xe2x80x94NR2C(O)CR2R2axe2x80x94, xe2x80x94CR2R2aC(O)NR2xe2x80x94, xe2x80x94CR2R2aNR2C(O)xe2x80x94, xe2x80x94NR2C(O)NR2xe2x80x94, xe2x80x94NR2xe2x80x94, xe2x80x94NR2CR2R2axe2x80x94, xe2x80x94CR2R2aNR2xe2x80x94, O, xe2x80x94CR2R2aOxe2x80x94, and xe2x80x94OCR2R2axe2x80x94;\nY is selected from one of the following carbocyclic and heterocyclic systems that are substituted with 0-2 R4a;\ncyclopropyl, cyclopentyl, cyclohexyl, phenyl, piperidinyl, piperazinyl, pyridyl, pyrimidyl, furanyl, morpholinyl, thiophenyl, pyrrolyl, pyrrolidinyl, oxazolyl, isoxazolyl, isoxazolinyl, thiazolyl, isothiazolyl, pyrazolyl, imidazolyl, oxadiazolyl, thiadiazolyl, triazolyl, 1,2,3-oxadiazolyl, 1,2,4-oxadiazolyl, 1,2,5-oxadiazolyl, 1,3,4-oxadiazolyl, 1,2,3-thiadiazolyl, 1,2,4-thiadiazolyl, 1,2,5-thiadiazolyl, 1,3,4-thiadiazolyl, 1,2,3-triazolyl, 1,2,4-triazolyl, 1,2,5-triazolyl, 1,3,4-triazolyl, benzofuranyl, benzothiofuranyl, indolyl, benzimidazolyl, benzoxazolyl, benzthiazolyl, indazolyl, benzisoxazolyl, benzisothiazolyl, and isoindazolyl;\nalternatively, Y is selected from the following bicyclic heteroaryl ring systems: \nK is selected from O, S, NH, and N.\nIn another embodiment, the present invention provides a novel compound wherein:\nY is selected from one of the following carbocyclic and heterocyclic systems which are substituted with 0-2 R4a;\nphenyl, piperidinyl, piperazinyl, pyridyl, pyrimidyl, furanyl, morpholinyl, thiophenyl, pyrrolyl, pyrrolidinyl, oxazolyl, isoxazolyl, thiazolyl, isothiazolyl, pyrazolyl, imidazolyl, oxadiazole, thiadiazole, triazole, 1,2,3-oxadiazole, 1,2,4-oxadiazole, 1,2,5-oxadiazole, 1,3,4-oxadiazole, 1,2,3-thiadiazole, 1,2,4-thiadiazole, 1,2,5-thiadiazole, 1,3,4-thiadiazole, 1,2,3-triazole, 1,2,4-triazole, 1,2,5-triazole, 1,3,4-triazole, benzofuran, benzothiofuran, indole, benzimidazole, benzimidazolone, benzoxazole, benzthiazole, indazole, benzisoxazole, benzisothiazole, and isoindazole.\nIn another embodiment, the present invention provides a novel compound wherein B is selected from phenyl, pyrrolidino, N-pyrrolidino-carbonyl, morpholino, N-morpholino-carbonyl, 1,2,3-triazolyl, imidazolyl, and benzimidazolyl, and is substituted with 0-1 R4a.\nIn another embodiment, the present invention provides a novel compound wherein B is selected from the group: 2-(aminosulfonyl)phenyl, 2-(methylaminosulfonyl)phenyl, 1-pyrrolidinocarbonyl, 2-(methylsulfonyl)phenyl, 2-(N,N-dimethylaminomethyl)phenyl, 2-(N-methylaminomethyl)phenyl, 2-(N-ethyl-N-methylaminomethyl)phenyl, 2-(N-pyrrolidinylmethyl)phenyl, 1-methyl-2-imidazolyl, 2-methyl-1-imidazolyl, 2-(dimethylaminomethyl)-1-imidazolyl, 2-(methylaminomethyl)-1-imidazolyl, 2-(N-(cyclopropylmethyl)aminomethyl)phenyl, 2-(N-(cyclobutyl)aminomethyl)phenyl, 2-(N-(cyclopentyl)aminomethyl)phenyl, 2-(N-(4-hydroxypiperidinyl)methyl)phenyl, and 2-(N-(3-hydroxypyrrolidinyl)methyl)phenyl.\nIn another embodiment, the present invention provides novel pharmaceutical compositions, comprising: a pharmaceutically acceptable carrier and a therapeutically effective amount of a compound of the present invention or a pharmaceutically acceptable salt form thereof.\nIn another embodiment, the present invention provides a novel method for treating or preventing a thromboembolic disorder, comprising: administering to a patient in need thereof a therapeutically effective amount of a compound of the present invention or a pharmaceutically acceptable salt form thereof.\nIn another embodiment, the present invention provides a compound of the present invention as described above for use in therapy.\nIn another embodiment, the present invention provides the use of the present invention as described above for the manufacture of a medicament for the treatment of a thromboembolic disorder.\nThe compounds herein described may have asymmetric centers. Compounds of the present invention containing an asymmetrically substituted atom may be isolated in optically active or racemic forms. It is well known in the art how to prepare optically active forms, such as by resolution of racemic forms or by synthesis from optically active starting materials. Many geometric isomers of olefins, Cxe2x95x90N double bonds, and the like can also be present in the compounds described herein, and all such stable isomers are contemplated in the present invention. Cis and trans geometric isomers of the compounds of the present invention are described and may be isolated as a mixture of isomers or as separated isomeric forms. All chiral, diastereomeric, racemic forms and all geometric isomeric forms of a structure are intended, unless the specific stereochemistry or isomeric form is specifically indicated. All processes used to prepare compounds of the present invention and intermediates made therein are considered to be part of the present invention. All tautomers of shown or described compounds are also considered to be part of the present invention.\nThe term xe2x80x9csubstituted,xe2x80x9d as used herein, means that any one or more hydrogens on the designated atom is replaced with a selection from the indicated group, provided that the designated atom\"\"s normal valency is not exceeded, and that the substitution results in a stable compound. When a substituent is keto (i.e., xe2x95x90O), then 2 hydrogens on the atom are replaced. Keto substituents are not present on aromatic moieties.\nThe present invention is intended to include all isotopes of atoms occurring in the present compounds. Isotopes include those atoms having the same atomic number but different mass numbers. By way of general example and without limitation, isotopes of hydrogen include tritium and deuterium. Isotopes of carbon include C-13 and C-14.\nWhen any variable (e.g., R6) occurs more than one time in any constituent or formula for a compound, its definition at each occurrence is independent of its definition at every other occurrence. Thus, for example, if a group is shown to be substituted with 0-2 R6, then said group may optionally be substituted with up to two R6 groups and R6 at each occurrence is selected independently from the definition of R6. Also, combinations of substituents and/or variables are permissible only if such combinations result in stable compounds.\nWhen a bond to a substituent is shown to cross a bond connecting two atoms in a ring, then such substituent may be bonded to any atom on the ring. When a substituent is listed without indicating the atom via which such substituent is bonded to the rest of the compound of a given formula, then such substituent may be bonded via any atom in such substituent. Combinations of substituents and/or variables are permissible only if such combinations result in stable compounds.\nAs used herein, xe2x80x9calkylxe2x80x9d is intended to include both branched and straight-chain saturated aliphatic hydrocarbon groups having the specified number of carbon atoms. C1-10 alkyl, is intended to include C1, C2, C3, C4, C5, C6, C7, C8, C9, and C10 alkyl groups. Examples of alkyl include, but are not limited to, methyl, ethyl, n-propyl, i-propyl, n-butyl, s-butyl, t-butyl, n-pentyl, and s-pentyl. xe2x80x9cHaloalkylxe2x80x9d is intended to include both branched and straight-chain saturated aliphatic hydrocarbon groups having the specified number of carbon atoms, substituted with 1 or more halogen (for example xe2x80x94CvFw where v=1 to 3 and w=1 to (2v+1)). Examples of haloalkyl include, but are not limited to, trifluoromethyl, trichloromethyl, pentafluoroethyl, and pentachloroethyl. xe2x80x9cAlkoxyxe2x80x9d represents an alkyl group as defined above with the indicated number of carbon atoms attached through an oxygen bridge. C1-10 alkoxy, is intended to include C1, C2, C3, C4, C5, C6, C7, C8, C9, and C10 alkoxy groups. Examples of alkoxy include, but are not limited to, methoxy, ethoxy, n-propoxy, i-propoxy, n-butoxy, s-butoxy, t-butoxy, n-pentoxy, and s-pentoxy. xe2x80x9cCycloalkylxe2x80x9d is intended to include saturated ring groups, such as cyclopropyl, cyclobutyl, or cyclopentyl. C3-7 cycloalkyl is intended to include C3, C4, C5, C6, and C7 cycloalkyl groups. Alkenylxe2x80x9d is intended to include hydrocarbon chains of either straight or branched configuration and one or more unsaturated carbonxe2x80x94carbon bonds that may occur in any stable point along the chain, such as ethenyl and propenyl. C2-10 alkenyl is intended to include C2, C3, C4, C5, C6, C7, C8, C9, and C10 alkenyl groups. xe2x80x9cAlkynylxe2x80x9d is intended to include hydrocarbon chains of either straight or branched configuration and one or more triple carbonxe2x80x94carbon bonds that may occur in any stable point along the chain, such as ethynyl and propynyl. C2-10 Alkynyl is intended to include C2, C3, C4, C5, C6, C7, C8, C9, and C10 alkynyl groups.\nxe2x80x9cHaloxe2x80x9d or xe2x80x9chalogenxe2x80x9d as used herein refers to fluoro, chloro, bromo, and iodo; and xe2x80x9ccounterionxe2x80x9d is used to represent a small, negatively charged species such as chloride, bromide, hydroxide, acetate, and sulfate.\nAs used herein, xe2x80x9ccarbocyclexe2x80x9d or xe2x80x9ccarbocyclic residuexe2x80x9d is intended to mean any stable 3, 4, 5, 6, or 7-membered monocyclic or bicyclic or 7, 8, 9, 10, 11, 12, or 13-membered bicyclic or tricyclic, any of which may be saturated, partially unsaturated, or aromatic. Examples of such carbocycles include, but are not limited to, cyclopropyl, cyclobutyl, cyclopentyl, cyclohexyl, cycloheptyl, adamantyl, cyclooctyl, [3.3.0]bicyclooctane, [4.3.0]bicyclononane, [4.4.0]bicyclodecane, [2.2.2]bicyclooctane, fluorenyl, phenyl, naphthyl, indanyl, adamantyl, and tetrahydronaphthyl.\nAs used herein, the term xe2x80x9cheterocyclexe2x80x9d or xe2x80x9cheterocyclic systemxe2x80x9d is intended to mean a stable 5, 6, or 7-membered monocyclic or bicyclic or 7, 8, 9, or 10-membered bicyclic heterocyclic ring which is saturated, partially unsaturated or unsaturated (aromatic), and which consists of carbon atoms and 1, 2, 3, or 4 heteroatoms independently selected from the group consisting of N, NH, O and S and including any bicyclic group in which any of the above-defined heterocyclic rings is fused to a benzene ring. The nitrogen and sulfur heteroatoms may optionally be oxidized. The heterocyclic ring may be attached to its pendant group at any heteroatom or carbon atom that results in a stable structure. The heterocyclic rings described herein may be substituted on carbon or on a nitrogen atom if the resulting compound is stable. A nitrogen in the heterocycle may optionally be quaternized. It is preferred that when the total number of S and O atoms in the heterocycle exceeds 1, then these heteroatoms are not adjacent to one another. It is preferred that the total number of S and O atoms in the heterocycle is not more than 1. As used herein, the term xe2x80x9caromatic heterocyclic systemxe2x80x9d or xe2x80x9cheteroarylxe2x80x9d is intended to mean a stable 5, 6, or 7-membered monocyclic or bicyclic or 7, 8, 9, or 10-membered bicyclic heterocyclic aromatic ring which consists of carbon atoms and 1, 2, 3, or 4 heteroatoms independently selected from the group consisting of N, NH, O and S. It is to be noted that total number of S and O atoms in the aromatic heterocycle is not more than 1.\nExamples of heterocycles include, but are not limited to, acridinyl, azocinyl, benzimidazolyl, benzofuranyl, benzothiofuranyl, benzothiophenyl, benzoxazolyl, benzoxazolinyl, benzthiazolyl, benztriazolyl, benztetrazolyl, benzisoxazolyl, benzisothiazolyl, benzimidazolinyl, carbazolyl, 4aH-carbazolyl, carbolinyl, chromanyl, chromenyl, cinnolinyl, decahydroquinolinyl, 2H, 6H-1, 5, 2-dithiazinyl, dihydrofuro[2,3-b]tetrahydrofuran, furanyl, furazanyl, imidazolidinyl, imidazolinyl, imidazolyl, 1H-indazolyl, indolenyl, indolinyl, indolizinyl, indolyl, 3H-indolyl, isatinoyl, isobenzofuranyl, isochromanyl, isoindazolyl, isoindolinyl, isoindolyl, isoquinolinyl, isothiazolyl, isoxazolyl, methylenedioxyphenyl, morpholinyl, naphthyridinyl, octahydroisoquinolinyl, oxadiazolyl, 1,2,3-oxadiazolyl, 1,2,4-oxadiazolyl, 1,2,5-oxadiazolyl, 1,3,4-oxadiazolyl, oxazolidinyl, oxazolyl, oxindolyl, pyrimidinyl, phenanthridinyl, phenanthrolinyl, phenazinyl, phenothiazinyl, phenoxathinyl, phenoxazinyl, phthalazinyl, piperazinyl, piperidinyl, piperidonyl, 4-piperidonyl, piperonyl, pteridinyl, purinyl, pyranyl, pyrazinyl, pyrazolidinyl, pyrazolinyl, pyrazolyl, pyridazinyl, pyridooxazole, pyridoimidazole, pyridothiazole, pyridinyl, pyridyl, pyrimidinyl, pyrrolidinyl, pyrrolinyl, 2H-pyrrolyl, pyrrolyl, quinazolinyl, quinolinyl, 4H-quinolizinyl, quinoxalinyl, quinuclidinyl, tetrahydrofuranyl, tetrahydroisoquinolinyl, tetrahydroquinolinyl, tetrazolyl, 6H-1,2,5-thiadiazinyl, 1,2,3-thiadiazolyl, 1,2,4-thiadiazolyl, 1,2,5-thiadiazolyl, 1,3,4-thiadiazolyl, thianthrenyl, thiazolyl, thienyl, thienothiazolyl, thienooxazolyl, thienoimidazolyl, thiophenyl, triazinyl, 1,2,3-triazolyl, 1,2,4-triazolyl, 1,2,5-triazolyl, 1,3,4-triazolyl, and xanthenyl. Also included are fused ring and spiro compounds containing, for example, the above heterocycles.\nThe phrase xe2x80x9cpharmaceutically acceptablexe2x80x9d is employed herein to refer to those compounds, materials, compositions, and/or dosage forms which are, within the scope of sound medical judgment, suitable for use in contact with the tissues of human beings and animals without excessive toxicity, irritation, allergic response, or other problem or complication, commensurate with a reasonable benefit/risk ratio.\nAs used herein, xe2x80x9cpharmaceutically acceptable saltsxe2x80x9d refer to derivatives of the disclosed compounds wherein the parent compound is modified by making acid or base salts thereof. Examples of pharmaceutically acceptable salts include, but are not limited to, mineral or organic acid salts of basic residues such as amines; alkali or organic salts of acidic residues such as carboxylic acids; and the like. The pharmaceutically acceptable salts include the conventional non-toxic salts or the quaternary ammonium salts of the parent compound formed, for example, from non-toxic inorganic or organic acids. For example, such conventional non-toxic salts include those derived from inorganic acids such as hydrochloric, hydrobromic, sulfuric, sulfamic, phosphoric, nitric and the like; and the salts prepared from organic acids such as acetic, propionic, succinic, glycolic, stearic, lactic, malic, tartaric, citric, ascorbic, pamoic, maleic, hydroxymaleic, phenylacetic, glutamic, benzoic, salicylic, sulfanilic, 2-acetoxybenzoic, fumaric, toluenesulfonic, methanesulfonic, ethane disulfonic, oxalic, isethionic, and the like.\nThe pharmaceutically acceptable salts of the present invention can be synthesized from the parent compound that contains a basic or acidic moiety by conventional chemical methods. Generally, such salts can be prepared by reacting the free acid or base forms of these compounds with a stoichiometric amount of the appropriate base or acid in water or in an organic solvent, or in a mixture of the two; generally, non-aqueous media like ether, ethyl acetate, ethanol, isopropanol, or acetonitrile are preferred. Lists of suitable salts are found in Remington\"\"s Pharmaceutical Sciences, 17th ed., Mack Publishing Company, Easton, Pa., 1985, p. 1418, the disclosure of which is hereby incorporated by reference.\nSince prodrugs are known to enhance numerous desirable qualities of pharmaceuticals (e.g., solubility, bioavailability, manufacturing, etc. . . . ) the compounds of the present invention may be delivered in prodrug form. Thus, the present invention is intended to cover prodrugs of the presently claimed compounds, methods of delivering the same and compositions containing the same. xe2x80x9cProdrugsxe2x80x9d are intended to include any covalently bonded carriers that release an active parent drug of the present invention in vivo when such prodrug is administered to a mammalian subject. Prodrugs the present invention are prepared by modifying functional groups present in the compound in such a way that the modifications are cleaved, either in routine manipulation or in vivo, to the parent compound. Prodrugs include compounds of the present invention wherein a hydroxy, amino, or sulfhydryl group is bonded to any group that, when the prodrug of the present invention is administered to a mammalian subject, it cleaves to form a free hydroxyl, free amino, or free sulfhydryl group, respectively. Examples of prodrugs include, but are not limited to, acetate, formate and benzoate derivatives of alcohol and amine functional groups in the compounds of the present invention.\nxe2x80x9cStable compoundxe2x80x9d and xe2x80x9cstable structurexe2x80x9d are meant to indicate a compound that is sufficiently robust to survive isolation to a useful degree of purity from a reaction mixture, and formulation into an efficacious therapeutic agent.\nxe2x80x9cSubstitutedxe2x80x9d is intended to indicate that one or more hydrogens on the atom indicated in the expression using xe2x80x9csubstitutedxe2x80x9d is replaced with a selection from the indicated group(s), provided that the indicated atom\"\"s normal valency is not exceeded, and that the substitution results in a stable compound. When a substituent is keto (i.e., xe2x95x90O) group, then 2 hydrogens on the atom are replaced.\nxe2x80x9cTherapeutically effective amountxe2x80x9d is intended to include an amount of a compound of the present invention or an amount of the combination of compounds claimed effective to inhibit factor Xa. The combination of compounds is preferably a synergistic combination. Synergy, as described, for example, by Chou and Talalay, Adv. Enzyme Regul. 1984, 22:27-55, occurs when the effect (in this case, inhibition of factor Xa) of the compounds when administered in combination is greater than the additive effect of the compounds when administered alone as a single agent. In general, a synergistic effect is most clearly demonstrated at sub-optimal concentrations of the compounds. Synergy can be in terms of lower cytotoxicity, increased antiviral effect, or some other beneficial effect of the combination compared with the individual components.\nThe compounds of the present invention can be prepared in a number of ways known to one skilled in the art of organic synthesis. The compounds of the present invention can be synthesized using the methods described below, together with synthetic methods known in the art of synthetic organic chemistry, or by variations thereon as appreciated by those skilled in the art. Preferred methods include, but are not limited to, those described below. The reactions are performed in a solvent appropriate to the reagents and materials employed and suitable for the transformations being effected. It will be understood by those skilled in the art of organic synthesis that the functionality present on the molecule should be consistent with the transformations proposed. This will sometimes require a judgment to modify the order of the synthetic steps or to select one particular process scheme over another in order to obtain a desired compound of the invention. It will also be recognized that another major consideration in the planning of any synthetic route in this field is the judicious choice of the protecting group used for protection of the reactive functional groups present in the compounds described in this invention. An authoritative account describing the many alternatives to the trained practitioner is Greene and Wuts (Protective Groups In Organic Synthesis, Wiley and Sons, 1991). All references cited herein are hereby incorporated in their entirety herein by reference. \nThe heterocyclic system shown in Scheme 1(wherein Gxe2x80x94D can be xe2x80x94CN or halo) can be prepared according to the schemes shown in PCT applications WO98/28269, WO98/28282, WO/9857934 and WO99/32454. In Scheme 1 and the schemes that follow, xe2x80x94Gxe2x80x94D represents rings D1, D2, or D3 attached to linker G or precursors to this moiety.\nCompounds of this invention wherein Gxe2x80x94D is a cyano group can be manipulated to afford thiadiazoles, oxadiazoles, aminooxadiazoles, pyrazoles, triazoles, and triazolones as outlined in Scheme 2.\nOther heterocycles contained in this invention can also be obtained via methods shown in Scheme 3\nVarious analogs of the compounds contained in Scheme 4 can be made following the methodologies outlined in Schemes 1-3 and WO01/05784.\nCompounds of this invention wherein ring D is a pyridyl, pyrimidyl, pyridazine, or other heterocyclic ring systems can also be prepared via the Suzuki and or the Stille coupling techniques from easily accessible bromo-phenyl precursors shown in Scheme 5.\nAdditionally, other compounds of this invention can also be prepared by those knowledgeable in the art following the methodology outlined in Scheme 6.\nThe compounds prepared via Scheme 6 can be used to prepare similar analogs with other five membered heterocyclic core systems or other linker G1. These can then be further elaborated via the methodologies shown in earlier schemes to obtain compounds of this invention.\nFor other carbon based heterocyclic or aryl core systems the above methodologies apply but can be further modified via the method outlined in Scheme 7\nThe synthesis shown in Scheme 7 can be adopted for other carbon based heterocycles or aryl rings. In addition to the above other compounds of this invention can also be prepared via the method outlined in Scheme 8.\nAlternatively, compounds of these inventions can also be prepared via the methodology outlined in Scheme 9.\nThe synthetic methodology applied in Scheme 9 can also be applied to all other heterocyclic or aryl systems. In all the cases and schemes provided the intermediates wherein X (or D-G-) is a halogen or cyano group these in turn can be manipulated to the heterocyclic systems contained in this invention.\nThe compounds of the present invention have a group xe2x80x9cA-Bxe2x80x9d attached to ring M directly or via Z. Preparations of some of the rings M and the xe2x80x9cA-Bxe2x80x9d moieties can follow the same methods described in WO97/23212, WO97/30971, WO97/38984, WO98/01428, WO98/06694, WO98/28269, WO98/28282, WO98/57934, WO98/57937, and WO98/57951, the contents of which are incorporated herein by reference\nThe compounds of this invention are useful as anticoagulants for the treatment or prevention of thromboembolic disorders in mammals. The term xe2x80x9cthromboembolic disordersxe2x80x9d as used herein includes arterial or venous cardiovascular or cerebrovascular thromboembolic disorders, including, for example, unstable angina, first or recurrent myocardial infarction, ischemic sudden death, transient ischemic attack, stroke, atherosclerosis, venous thrombosis, deep vein thrombosis, thrombophlebitis, arterial embolism, coronary and cerebral arterial thrombosis, cerebral embolism, kidney embolisms, and pulmonary embolisms. The anticoagulant effect of compounds of the present invention is believed to be due to inhibition of factor Xa or thrombin.\nThe effectiveness of compounds of the present invention as inhibitors of factor Xa was determined using purified human factor Xa and synthetic substrate. The rate of factor Xa hydrolysis of chromogenic substrate S2222 (Kabi Pharmacia, Franklin, Ohio) was measured both in the absence and presence of compounds of the present invention. Hydrolysis of the substrate resulted in the release of pNA, which was monitored spectrophotometrically by measuring the increase in absorbance at 405 nM. A decrease in the rate of absorbance change at 405 nm in the presence of inhibitor is indicative of enzyme inhibition. The results of this assay are expressed as inhibitory constant, Ki.\nFactor Xa determinations were made in 0.10 M sodium phosphate buffer, pH 7.5, containing 0.20 M NaCl, and 0.5% PEG 8000. The Michaelis constant, Km, for substrate hydrolysis was determined at 25xc2x0 C. using the method of Lineweaver and Burk. Values of Ki were determined by allowing 0.2-0.5 nM human factor Xa (Enzyme Research Laboratories, South Bend, Ind.) to react with the substrate (0.20 mM-1 mM) in the presence of inhibitor. Reactions were allowed to go for 30 minutes and the velocities (rate of absorbance change vs time) were measured in the time frame of 25-30 minutes. The following relationship was used to calculate Ki values:\n(voxe2x88x92vs)/vs=I/(Ki (1+S/Km))\nwhere:\nvo is the velocity of the control in the absence of inhibitor;\nvs is the velocity in the presence of inhibitor;\nI is the concentration of inhibitor;\nKi is the dissociation constant of the enzyme:inhibitor complex;\nS is the concentration of substrate;\nKm is the Michaelis constant.\nUsing the methodology described above, some compounds of the present invention were found to exhibit a Ki of xe2x89xa610 xcexcM, thereby confirming the utility of the compounds of the present invention as effective Xa inhibitors.\nCompounds tested in the above assay are considered to be active if they exhibit a Ki of xe2x89xa610 xcexcM. Preferred compounds of the present invention have Ki\"\"s of xe2x89xa61 xcexcM. More preferred compounds of the present invention have Ki\"\"s of xe2x89xa60.1 xcexcM. Even more preferred compounds of the present invention have Ki\"\"s of xe2x89xa60.01 xcexcM. Still more preferred compounds of the present invention have Ki\"\"s of xe2x89xa60.001 xcexcM.\nThe antithrombotic effect of compounds of the present invention can be demonstrated in a rabbit arterio-venous (AV) shunt thrombosis model. In this model, rabbits weighing 2-3 kg anesthetized with a mixture of xylazine (10 mg/kg i.m.) and ketamine (50 mg/kg i.m.) are used. A saline-filled AV shunt device is connected between the femoral arterial and the femoral venous cannulae. The AV shunt device consists of a piece of 6-cm tygon tubing that contains a piece of silk thread. Blood will flow from the femoral artery via the AV-shunt into the femoral vein. The exposure of flowing blood to a silk thread will induce the formation of a significant thrombus. After forty minutes, the shunt is disconnected and the silk thread covered with thrombus is weighed. Test agents or vehicle will be given (i.v., i.p., s.c., or orally) prior to the opening of the AV shunt. The percentage inhibition of thrombus formation is determined for each treatment group. The ID50 values (dose which produces 50% inhibition of thrombus formation) are estimated by linear regression.\nThe compounds of the present invention may also be useful as inhibitors of serine proteases, notably human thrombin, plasma kallikrein and plasmin. Because of their inhibitory action, these compounds are indicated for use in the prevention or treatment of physiological reactions, blood coagulation and inflammation, catalyzed by the aforesaid class of enzymes. Specifically, the compounds have utility as drugs for the treatment of diseases arising from elevated thrombin activity such as myocardial infarction, and as reagents used as anticoagulants in the processing of blood to plasma for diagnostic and other commercial purposes.\nSome compounds of the present invention were shown to be direct acting inhibitors of the serine protease thrombin by their ability to inhibit the cleavage of small molecule substrates by thrombin in a purified system. In vitro inhibition constants were determined by the method described by Kettner et al. in J. Biol. Chem. 1990, 265, 18289-18297, herein incorporated by reference. In these assays, thrombin-mediated hydrolysis of the chromogenic substrate S2238 (Helena Laboratories, Beaumont, Tex.) was monitored spectrophotometrically. Addition of an inhibitor to the assay mixture results in decreased absorbance and is indicative of thrombin inhibition. Human thrombin (Enzyme Research Laboratories, Inc., South Bend, Ind.) at a concentration of 0.2 nM in 0.10 M sodium phosphate buffer, pH 7.5, 0.20 M NaCl, and 0.5% PEG 6000, was incubated with various substrate concentrations ranging from 0.20 to 0.02 mM. After 25 to 30 minutes of incubation, thrombin activity was assayed by monitoring the rate of increase in absorbance at 405 nm that arises owing to substrate hydrolysis. Inhibition constants were derived from reciprocal plots of the reaction velocity as a function of substrate concentration using the standard method of Lineweaver and Burk. Using the methodology described above, some compounds of this invention were evaluated and found to exhibit a Ki of less than 10 xcexcm, thereby confirming the utility of the compounds of the present invention as effective thrombin inhibitors.\nThe compounds of the present invention can be administered alone or in combination with one or more additional therapeutic agents. These include other anti-coagulant or coagulation inhibitory agents, anti-platelet or platelet inhibitory agents, thrombin inhibitors, or thrombolytic or fibrinolytic agents.\nThe compounds are administered to a mammal in a therapeutically effective amount. By xe2x80x9ctherapeutically effective amountxe2x80x9d it is meant an amount of a compound of the present invention that, when administered alone or in combination with an additional therapeutic agent to a mammal, is effective to prevent or ameliorate the thromboembolic disease condition or the progression of the disease.\nBy xe2x80x9cadministered in combinationxe2x80x9d or xe2x80x9ccombination therapyxe2x80x9d it is meant that the compound of the present invention and one or more additional therapeutic agents are administered concurrently to the mammal being treated. When administered in combination each component may be administered at the same time or sequentially in any order at different points in time. Thus, each component may be administered separately but sufficiently closely in time so as to provide the desired therapeutic effect. Other anticoagulant agents (or coagulation inhibitory agents) that may be used in combination with the compounds of this invention include warfarin and heparin, as well as other factor Xa inhibitors such as those described in the publications identified above under Background of the Invention.\nThe term anti-platelet agents (or platelet inhibitory agents), as used herein, denotes agents that inhibit platelet function such as by inhibiting the aggregation, adhesion or granular secretion of platelets. Such agents include, but are not limited to, the various known non-steroidal anti-inflammatory drugs (NSAIDS) such as aspirin, ibuprofen, naproxen, sulindac, indomethacin, mefenamate, droxicam, diclofenac, sulfinpyrazone, and piroxicam, including pharmaceutically acceptable salts or prodrugs thereof. Of the NSAIDS, aspirin (acetylsalicyclic acid or ASA), and piroxicam are preferred. Other suitable anti-platelet agents include ticlopidine, including pharmaceutically acceptable salts or prodrugs thereof. Ticlopidine is also a preferred compound since it is known to be gentle on the gastro-intestinal tract in use. Still other suitable platelet inhibitory agents include IIb/IIIa antagonists, thromboxane-A2-receptor antagonists and thromboxane-A2-synthetase inhibitors, as well as pharmaceutically acceptable salts or prodrugs thereof.\nThe term thrombin inhibitors (or anti-thrombin agents), as used herein, denotes inhibitors of the serine protease thrombin. By inhibiting thrombin, various thrombin-mediated processes, such as thrombin-mediated platelet activation (that is, for example, the aggregation of platelets, and/or the granular secretion of plasminogen activator inhibitor-1 and/or serotonin) and/or fibrin formation are disrupted. A number of thrombin inhibitors are known to one of skill in the art and these inhibitors are contemplated to be used in combination with the present compounds. Such inhibitors include, but are not limited to, boroarginine derivatives, boropeptides, heparins, hirudin and argatroban, including pharmaceutically acceptable salts and prodrugs thereof. Boroarginine derivatives and boropeptides include N-acetyl and peptide derivatives of boronic acid, such as C-terminal a-aminoboronic acid derivatives of lysine, ornithine, arginine, homoarginine and corresponding isothiouronium analogs thereof. The term hirudin, as used herein, includes suitable derivatives or analogs of hirudin, referred to herein as hirulogs, such as disulfatohirudin. Boropeptide thrombin inhibitors include compounds described in Kettner et al., U.S. Pat. No. 5,187,157 and European Patent Application Publication No. 293 881 A2, the disclosures of which are hereby incorporated herein by reference. Other suitable boroarginine derivatives and boropeptide thrombin inhibitors include those disclosed in PCT Application Publication No. 92/07869 and European Patent Application Publication No. 471,651 A2, the disclosures of which are hereby incorporated herein by reference.\nThe term thrombolytics (or fibrinolytic) agents (or thrombolytics or fibrinolytics), as used herein, denotes agents that lyse blood clots (thrombi). Such agents include tissue plasminogen activator, anistreplase, urokinase or streptokinase, including pharmaceutically acceptable salts or prodrugs thereof. The term anistreplase, as used herein, refers to anisoylated plasminogen streptokinase activator complex, as described, for example, in European Patent Application No. 028,489, the disclosure of which is hereby incorporated herein by reference herein. The term urokinase, as used herein, is intended to denote both dual and single chain urokinase, the latter also being referred to herein as prourokinase.\nAdministration of the compounds of the present invention in combination with such additional therapeutic agent, may afford an efficacy advantage over the compounds and agents alone, and may do so while permitting the use of lower doses of each. A lower dosage minimizes the potential of side effects, thereby providing an increased margin of safety.\nThe compounds of the present invention are also useful as standard or reference compounds, for example as a quality standard or control, in tests or assays involving the inhibition of factor Xa. Such compounds may be provided in a commercial kit, for example, for use in pharmaceutical research involving factor Xa. For example, a compound of the present invention could be used as a reference in an assay to compare its known activity to a compound with an unknown activity. This would ensure the experimenter that the assay was being performed properly and provide a basis for comparison, especially if the test compound was a derivative of the reference compound. When developing new assays or protocols, compounds according to the present invention could be used to test their effectiveness.\nThe compounds of the present invention may also be used in diagnostic assays involving factor Xa. For example, the presence of factor Xa in an unknown sample could be determined by addition of chromogenic substrate S2222 to a series of solutions containing test sample and optionally one of the compounds of the present invention. If production of pNA is observed in the solutions containing test sample, but not in the presence of a compound of the present invention, then one would conclude factor Xa was present.\nThe compounds of this invention can be administered in such oral dosage forms as tablets, capsules (each of which includes sustained release or timed release formulations), pills, powders, granules, elixirs, tinctures, suspensions, syrups, and emulsions. They may also be administered in intravenous (bolus or infusion), intraperitoneal, subcutaneous, or intramuscular form, all using dosage forms well known to those of ordinary skill in the pharmaceutical arts. They can be administered alone, but generally will be administered with a pharmaceutical carrier selected on the basis of the chosen route of administration and standard pharmaceutical practice.\nThe dosage regimen for the compounds of the present invention will, of course, vary depending upon known factors, such as the pharmacodynamic characteristics of the particular agent and its mode and route of administration; the species, age, sex, health, medical condition, and weight of the recipient; the nature and extent of the symptoms; the kind of concurrent treatment; the frequency of treatment; the route of administration, the renal and hepatic function of the patient, and the effect desired. A physician or veterinarian can determine and prescribe the effective amount of the drug required to prevent, counter, or arrest the progress of the thromboembolic disorder.\nBy way of general guidance, the daily oral dosage of each active ingredient, when used for the indicated effects, will range between about 0.001 to 1000 mg/kg of body weight, preferably between about 0.01 to 100 mg/kg of body weight per day, and most preferably between about 1.0 to 20 mg/kg/day. Intravenously, the most preferred doses will range from about 1 to about 10 mg/kg/minute during a constant rate infusion. Compounds of this invention may be administered in a single daily dose, or the total daily dosage may be administered in divided doses of two, three, or four times daily.\nCompounds of this invention can be administered in intranasal form via topical use of suitable intranasal vehicles, or via transdermal routes, using transdermal skin patches. When administered in the form of a transdermal delivery system, the dosage administration will, of course, be continuous rather than intermittent throughout the dosage regimen.\nThe compounds are typically administered in admixture with suitable pharmaceutical diluents, excipients, or carriers (collectively referred to herein as pharmaceutical carriers) suitably selected with respect to the intended form of administration, that is, oral tablets, capsules, elixirs, syrups and the like, and consistent with conventional pharmaceutical practices.\nFor instance, for oral administration in the form of a tablet or capsule, the active drug component can be combined with an oral, non-toxic, pharmaceutically acceptable, inert carrier such as lactose, starch, sucrose, glucose, methyl callulose, magnesium stearate, dicalcium phosphate, calcium sulfate, mannitol, sorbitol and the like; for oral administration in liquid form, the oral drug components can be combined with any oral, non-toxic, pharmaceutically acceptable inert carrier such as ethanol, glycerol, water, and the like. Moreover, when desired or necessary, suitable binders, lubricants, disintegrating agents, and coloring agents can also be incorporated into the mixture. Suitable binders include starch, gelatin, natural sugars such as glucose or beta-lactose, corn sweeteners, natural and synthetic gums such as acacia, tragacanth, or sodium alginate, carboxymethylcellulose, polyethylene glycol, waxes, and the like. Lubricants used in these dosage forms include sodium oleate, sodium stearate, magnesium stearate, sodium benzoate, sodium acetate, sodium chloride, and the like. Disintegrators include, without limitation, starch, methyl cellulose, agar, bentonite, xanthan gum, and the like.\nThe compounds of the present invention can also be administered in the form of liposome delivery systems, such as small unilamellar vesicles, large unilamellar vesicles, and multilamellar vesicles. Liposomes can be formed from a variety of phospholipids, such as cholesterol, stearylamine, or phosphatidylcholines.\nCompounds of the present invention may also be coupled with soluble polymers as targetable drug carriers. Such polymers can include polyvinylpyrrolidone, pyran copolymer, polyhydroxypropylmethacrylamide-phenol, polyhydroxyethylaspartamidephenol, or polyethyleneoxidepolylysine substituted with palmitoyl residues. Furthermore, the compounds of the present invention may be coupled to a class of biodegradable polymers useful in achieving controlled release of a drug, for example, polylactic acid, polyglycolic acid, copolymers of polylactic and polyglycolic acid, polyepsilon caprolactone, polyhydroxy butyric acid, polyorthoesters, polyacetals, polydihydropyrans, polycyanoacylates, and crosslinked or amphipathic block copolymers of hydrogels.\nDosage forms (pharmaceutical compositions) suitable for administration may contain from about 1 milligram to about 100 milligrams of active ingredient per dosage unit. In these pharmaceutical compositions the active ingredient will ordinarily be present in an amount of about 0.5-95% by weight based on the total weight of the composition.\nGelatin capsules may contain the active ingredient and powdered carriers, such as lactose, starch, cellulose derivatives, magnesium stearate, stearic acid, and the like. Similar diluents can be used to make compressed tablets. Both tablets and capsules can be manufactured as sustained release products to provide for continuous release of medication over a period of hours. Compressed tablets can be sugar coated or film coated to mask any unpleasant taste and protect the tablet from the atmosphere, or enteric coated for selective disintegration in the gastrointestinal tract.\nLiquid dosage forms for oral administration can contain coloring and flavoring to increase patient acceptance.\nIn general, water, a suitable oil, saline, aqueous dextrose (glucose), and related sugar solutions and glycols such as propylene glycol or polyethylene glycols are suitable carriers for parenteral solutions. Solutions for parenteral administration preferably contain a water soluble salt of the active ingredient, suitable stabilizing agents, and if necessary, buffer substances. Antioxidizing agents such as sodium bisulfite, sodium sulfite, or ascorbic acid, either alone or combined, are suitable stabilizing agents. Also used are citric acid and its salts and sodium EDTA. In addition, parenteral solutions can contain preservatives, such as benzalkonium chloride, methyl- or propyl-paraben, and chlorobutanol.\nSuitable pharmaceutical carriers are described in Remington\"\"s Pharmaceutical Sciences, Mack Publishing Company, a standard reference text in this field.\nRepresentative useful pharmaceutical dosage-forms for administration of the compounds of this invention can be illustrated as follows:\nCapsules\nA large number of unit capsules can be prepared by filling standard two-piece hard gelatin capsules each with 100 milligrams of powdered active ingredient, 150 milligrams of lactose, 50 milligrams of cellulose, and 6 milligrams magnesium stearate.\nSoft Gelatin Capsules\nA mixture of active ingredient in a digestable oil such as soybean oil, cottonseed oil or olive oil may be prepared and injected by means of a positive displacement pump into gelatin to form soft gelatin capsules containing 100 milligrams of the active ingredient. The capsules should be washed and dried.\nTablets\nTablets may be prepared by conventional procedures so that the dosage unit is 100 milligrams of active ingredient, 0.2 milligrams of colloidal silicon dioxide, 5 milligrams of magnesium stearate, 275 milligrams of microcrystalline cellulose, 11 milligrams of starch and 98.8 milligrams of lactose. Appropriate coatings may be applied to increase palatability or delay absorption.\nInjectable\nA parenteral composition suitable for administration by injection may be prepared by stirring 1.5% by weight of active ingredient in 10% by volume propylene glycol and water. The solution should be made isotonic with sodium chloride and sterilized.\nSuspension\nAn aqueous suspension can be prepared for oral administration so that each 5 mL contain 100 mg of finely divided active ingredient, 200 mg of sodium carboxymethyl cellulose, 5 mg of sodium benzoate, 1.0 g of sorbitol solution, U.S.P., and 0.025 mL of vanillin.\nWhere the compounds of this invention are combined with other anticoagulant agents, for example, a daily dosage may be about 0.1 to 100 milligrams of the compound of the present invention and about 1 to 7.5 milligrams of the second anticoagulant, per kilogram of patient body weight. For a tablet dosage form, the compounds of this invention generally may be present in an amount of about 5 to 10 milligrams per dosage unit, and the second anti-coagulant in an amount of about 1 to 5 milligrams per dosage unit.\nWhere the compounds of the present invention are administered in combination with an anti-platelet agent, by way of general guidance, typically a daily dosage may be about 0.01 to 25 milligrams of the compound of the present invention and about 50 to 150 milligrams of the anti-platelet agent, preferably about 0.1 to 1 milligrams of the compound of the present invention and about 1 to 3 milligrams of anti-platelet agents, per kilogram of patient body weight.\nWhere the compounds of the present invention are administered in combination with thrombolytic agent, typically a daily dosage may be about 0.1 to 1 milligrams of the compound of the present invention, per kilogram of patient body weight and, in the case of the thrombolytic agents, the usual dosage of the thrombolyic agent when administered alone may be reduced by about 70-80% when administered with a compound of the present invention.\nWhere two or more of the foregoing second therapeutic agents are administered with the compound of the present invention, generally the amount of each component in a typical daily dosage and typical dosage form may be reduced relative to the usual dosage of the agent when administered alone, in view of the additive or synergistic effect of the therapeutic agents when administered in combination.\nParticularly when provided as a single dosage unit, the potential exists for a chemical interaction between the combined active ingredients. For this reason, when the compound of the present invention and a second therapeutic agent are combined in a single dosage unit they are formulated such that although the active ingredients are combined in a single dosage unit, the physical contact between the active ingredients is minimized (that is, reduced). For example, one active ingredient may be enteric coated. By enteric coating one of the active ingredients, it is possible not only to minimize the contact between the combined active ingredients, but also, it is possible to control the release of one of these components in the gastrointestinal tract such that one of these components is not released in the stomach but rather is released in the intestines. One of the active ingredients may also be coated with a material that effects a sustained-release throughout the gastrointestinal tract and also serves to minimize physical contact between the combined active ingredients. Furthermore, the sustained-released component can be additionally enteric coated such that the release of this component occurs only in the intestine. Still another approach would involve the formulation of a combination product in which the one component is coated with a sustained and/or enteric release polymer, and the other component is also coated with a polymer such as a low-viscosity grade of hydroxypropyl methylcellulose (HPMC) or other appropriate materials as known in the art, in order to further separate the active components. The polymer coating serves to form an additional barrier to interaction with the other component.\nThese as well as other ways of minimizing contact between the components of combination products of the present invention, whether administered in a single dosage form or administered in separate forms but at the same time by the same manner, will be readily apparent to those skilled in the art, once armed with the present disclosure.\nOther features of the invention will become apparent in the course of the following descriptions of exemplary embodiments that are given for illustration of the invention and are not intended to be limiting thereof."}
{"text": "The structure of the growth hormone release inhibiting factor, somatostatin, has been determined by Brazeau et al., Science, 179, 77(1973). Several techniques for synthesizing somatostatin have been reported in the literature, including the solid phase method of Rivier, J.A.C.S. 96, 2986(1974) and the solution methods of Sarantakis et al., Biochemical Biophysical Research Communications 54, 234(1973) and Immer et al., Helv. Chim. Acta, 57, 730(1974)."}
{"text": "1. Field of Application and Prior Art\nThe invention relates to an appliance for switching on and off several heating devices of a cooker, as well as a cooker having such an appliance.\nRadiant heaters with a diameter which can exceed 230 mm exist for hobs with glass ceramic plates, for example. They in part suffer from the problem that an energy or power supply through so-called energy or power control devices on the one hand and an excess temperature protection for the glass ceramic plate over the radiant heater through so-called temperature limiters on the other are limited by the maximum power levels which can be applied and by a so-called flicker standard. The flicker standard indicates how frequently in a specific time period a specific power may be switched on and off for a cooker and is intended to prevent significant supply network reactive effects in line with the power supply companies. The switching capacity of both the power control devices and the temperature limiters, which operate with so-called snap-action switches and such as are for example described in U.S. Pat. Nos. 6,064,045 and 4,633,238, is generally limited. In the USA this is 12 or 13 Ampere, for example, and consequently 100,000 switching cycles must be attainable.\nThus, with the normally predetermined mains voltage, it is not possible to further increase the power of a radiant heater.\n2. Problem and Solution\nThe problem of the invention is to provide the aforementioned appliance and aforementioned cooker enabling the prior art problems to be avoided and which in particular enables the maximum power to be increased, particularly for a hotplate with radiant heating.\nThis problem is solved by an appliance having the features of claim 1 and a cooker having the features of claim 11. Advantageous and preferred developments of the invention form the subject matter of the further claims and are explained in greater detail hereinafter. By express reference the wording of the claims is made into part of the content of the description. In line with the present application “have” means that this feature can be inter alia provided, independently of further features.\nAccording to the invention, the appliance has a temperature detection device or in certain circumstances forms the latter and has a temperature-dependent, thermomechanical release or tripping movement as the operating principle. The temperature detection device has a first switching device or is connected thereto and can be directly operated at an adjustable tripping point by a tripping movement in order to switch on and off the first switching device. A second switching device is provided for switching on and off the second heating device. Said second switching device is also activatable by the tripping movement. Both the first and second switching devices are designed for switching the heating power of the particular heating device, that means as so-called power switches, relays, for example. They are so constructed that they achieve the necessary number of switching cycles. For the second switching device or the tripping thereof, it can be directly or indirectly activated by the tripping movement. In the case of a direct activation the switching contacts are directly moved by the tripping movement. With indirect activation in the sense of the present application the tripping of the first switching device is activated by the tripping movement and the signal caused by the same or a switched voltage activates the second switching device. Thus, here activation does not take place in direct mechanical manner via the tripping movement.\nThus, it is possible as a result of the invention to create a temperature detection device which, in place of a conventional temperature limiter, such as can be gathered from U.S. Pat. No. 4,633,238, for example, can be used for monitoring the glass ceramic plate of a hob with radiant heater with respect to excess temperatures. The two heating devices belong to a single hotplate, that means at least on occasions are operated jointly. In particular there is a two or multiple circuit hotplate. Thus, the temperature detection device trips the first and second switching devices and each switching device switches on or off the particular associated heating device. Thus, in the case of an excess temperature, the total power of the hotplate no longer has to be switched by a single switching device of the temperature detection device, so that the total power of the heating devices or hotplate can be increased.\nPreferably, when using relays as switching devices, they are so designed that when a current is not flowing through them, that means in the deactivated or uncontrolled state, they are opened. The heating devices connected thereto are then switched off.\nAccording to a first possibility provided by the invention, the second switching device can be directly activated by the first switching device. This can for example take place in that with the first switching device closed for switching on the first heating device the second switching device receives a corresponding electric signal so as to then activate or switch on the second heating device. Alternatively and as a result of the connection of the first heating device to a supply voltage, the second switching device can be supplied with energy for activation. Thus, in this embodiment, the second switching device is not directly activated by the tripping movement and is instead indirectly activated, in that it is controlled by the first switching device activated by the tripping movement.\nIn another possibility provided by the invention, it is possible to provide a third switching device, which is connected or coupled to the temperature detection device and at an adjustable tripping point can be directly operated by the tripping movement. An operation advantageously takes place simultaneously with the first switching device, that means at the same tripping point. Such a third switching device can directly activate or supply energy to the second switching device, so that it correspondingly switches on or off the second heating device. Thus, here the second switching device is activated by the third switching device, unlike in the previous embodiment where this takes place through the first switching device. This offers the advantage that in this case no further functions or connections have to be provided for the first switching device, but it is necessary to provide the third switching device. However, it is also possible to use a switching device for this purpose which is already present in conventional temperature limiters, such as for example according to U.S. Pat. No. 4,633,238, for indicating a temperature of above 100° C., for example, so as to constitute a hot indicator.\nAccording to another possibility of the invention it is also possible to incorporate the second switching device into the temperature detection device in the same way as the first switching device and in particular with the same construction. In this case increased constructional demands are made, particularly with respect to insulation resistances and gaps.\nThe third switching device for controlling the second switching device can either operate with a lower voltage than the mains voltage or can be designed for significantly lower currents. Its function is merely to switch the activation current for the second switching device. In the case of power relays and in certain circumstances electronic switches this is relatively low, so that no increased mechanical or contact-specific demands have to be made on the third switching device.\nThe temperature detection device can advantageously have an expansion device, which expands or brings about a mechanical tripping movement as a function of the detected temperature. It can be a temperature limiter according to U.S. Pat. No. 4,633,238, to which express reference is made. Such a temperature limiter can have two elongated, parallel longitudinal elements, which are coupled at one end and have different temperature expansion coefficients. At a predetermined or detected temperature they expand to a varying degree and this relative movement leads to the tripping movement, because the relative movement between them occurs at the uncoupled end. For this purpose one of the elements can be fixed to the temperature detection device. The other moves relative thereto and gives rise to a force which, together with the tripping movement, can for example operate a switching device or the like.\nAccording to a further development the appliance can have a control for the heating devices or can be connected thereto, for example in the form of a so-called timing or cyclic energy control device. As a result thereof radiant heaters, for example, can be operated cyclically, which means that they are either switched off or switched on at full power. Therefore the level of the energy generation at the heating devices can be predetermined over several cycles or over and beyond a specific time. This is described in U.S. Pat. No. 6,064,045 to which express reference is made. When the control is switched off none of the aforementioned switching devices is activated or activatable. This ensures that the level of the energy generation or the nature of the timing is predetermined by the control or the energy control device. The temperature detection device advantageously merely provides an overheating protection of a glass ceramic plate, for example, against excessive heating by the heating devices.\nA cooker according to the invention can contain an aforementioned appliance for at least two independently controllable heating devices. These two heating devices form a hotplate or are very closely juxtaposed, so that in some modes they can be operated jointly for heating a single correspondingly large, stood-on cooking vessel. The first heating device is the main heating device, whereas the second heating device is an additional heating device. The second device can either be connected laterally and roughly over roughly a half to the first heating device and can optionally even completely surround the same. For smaller cooking vessels, they are heated solely by the first heating device. With larger cooking vessels covering the entire surface area of the first and second heating devices, both heating devices are used jointly. With such a joint use, it is naturally advantageous if the temperature detection device in the case of an excess temperature in an area, particularly of the first heating device, not only switches off the latter, but in fact both heating devices. Otherwise and in particular when using radiant heaters, an unaccustomed picture would arise for the user.\nAdvantageously the power of the first heating device is significantly higher than that of the second heating device and is roughly twice as high, for example.\nFor the second heating device it is possible to provide a maximum continuous power output corresponding to a power density of max. approximately 2.5 W/cmý. Such a value has the advantage that when using radiant heaters and a glass ceramic plate as the hob, such a power density can be operated so-to-speak in unprotected form, that is without any temperature monitoring for the glass ceramic. Dangerous temperatures of approximately 550° C. and higher cannot be obtained. Such power densities cannot lead to an overheating of the glass ceramic. The power density can also be selected above 2.5 W/cmý if the glass ceramic manufacture's specifications or tests reveal or allow this.\nThis makes it possible for the temperature device in its function as an overheating protection to only have to monitor the first heating device, but not the second heating device. It can for example cover the first heating device, but not the second heating device. It can alternatively be deactivated in the area of the second heating device. For this purpose it is possible, for example in the case of a rod control unit in the area to be deactivated of one of the two longitudinal elements to be partly replaced by another partial longitudinal element, which has the same temperature expansion as the other longitudinal element. Thus, by heating in the area to be deactivated, there is no relative movement or no contribution to a relative movement.\nThe cooker can advantageously have a timing or cyclic energy control device, as described hereinbefore and as can for example be gathered from U.S. Pat. No. 6,064,045. This energy control device is constructed for at least for activating the first heating device. Due to the fact that the second heating device is controlled or activated by means of the aforementioned appliance or temperature detection device and the power level of the heating devices is determined by the energy control device, this predetermines the on and off times both for the first heating device and consequently the second heating device.\nAccording to a further development a fourth switching device can be provided and which, independently of the second or third switching device, activates or deactivates the second heating device for synchronous operation with the heating device. This fourth switching device can for example make it possible for a user, as a function of the size of the stood-on cooking vessel, to activate or not activate the additional heater. It can for example be formed by a power switch looped into the supply voltage and in particular a relay, with a corresponding control for a user, for example a contact switch. The second switching device can be provided separately from the energy control device on the cooker. In the manner described hereinbefore, it can be a separate power switch, particularly a relay.\nThese and further features can be gathered from the claims, the description and the drawings and the individual features, both singly or in the form of subcombinations, can be implemented in an embodiment of the invention and in other fields and can represent advantageous, independently protectable constructions for which protection is claimed here. The subdivision of the application into individual sections and the subheadings in no way limit the general validity of the statements made thereunder."}
{"text": "Inorganic phosphate (P.sub.i), a charged anion, is essential to bioenergetics, metabolic regulation, and bone and membrane structure. It is well known that P.sub.i homeostasis in the body depends primarily on mechanisms that govern the renal excretion of P.sub.i into the glomerular filtrate and its subsequent reabsorption against an electrochemical gradient via brush-border epithelial cells located in the proximal tubule of the kidney [J. Bonjour and J. Caverzasio, Reviews in Physiological Pharmacology, 100:161-214 (1985); V. W. Dennis, Phosphate homeostasis, in HANDBOOK OF PHYSIOLOGY, (S. Shultz, ed. 1991) at pages 1785-1815.] This transepithelial transport of P.sub.i is mediated, in part, by a transport system which is driven by the transmembrane Na.sup.+ gradient across the microvilli brush border membrane. However, it remains largely unknown how cells transport and regulate necessary the intracellular concentrations of P.sub.i, and the molecular events underlying this system. Experiments using isolated kidney tubules or brush-border membranes have shown that P.sub.i transport is rather complex, regulated not only by extracellular [P.sub.i ] but also by neurotransmitters such as catecholamines (for review see V. W. Dennis, supra), and by a variety of hormones and metabolic factors. Berndt and Knox, \"Renal Regulation of Phosphate Excretion\", in, THE KIDNEY. PHYSIOLOGY AND PATHOPHYSIOLOGY, (D. W. Seldin and G. Giebisch, eds., 1991) at pages 1381-1396. Renal denervation, for example, decreases sodium and phosphate reabsorption. Norepinephrine released from nerve endings in proximity to renal tubules acts on the proximal tubule to increase phosphate reabsorption. In studies of isolated tubules, however, dopamine is shown to inhibit phosphate and sodium transport in the rabbit proximal tubule. Furthermore, several studies also show than depletion of extracellular P.sub.i or increased circulating levels of parathyroid hormone alter the activity and expression of transporter molecules or both.\nSeveral recent reports have demonstrated that P.sub.i homeostasis significantly affects the central nervous system (CNS). Phosphate/calcium alterations in serum, for example, have been implicated in the etiology and pathogenesis of Alzheimer's diseases. Depletion of high energy phosphates (phosphocreatine) and ATP is thought to be part of the final common pathway mediating excitotoxic neuronal cell death secondary to a wide variety of insults. Tight coupling between P.sub.i transport and ATP production has been observed in many cells and tissues. Chronic P.sub.i depletion in vivo is associated with a significant reduction in the ATP content of polymorphonuclear leukocytes, platelets, and various tissues including kidney, heart, and skeletal muscle. A similar observation has been made in cultured peripheral vagal nerves. This reduction in intracellular ATP has been shown to be a direct consequence of the decrease in intracellular P.sub.i which occurs following P.sub.i depletion. In addition to its possible role in ATP biosynthesis, several lines of evidence have suggested that P.sub.i may be involved in neuronal signalling events. In this regard, a study using brain tissue has recently shown that physiological concentrations of P.sub.i can enhance the ATP-dependent binding of Ca.sup.++ to brain microsomes, resulting in a larger intracellular pool of Ca.sup.++ releasable by inositol triphosphate. Our recent work have demonstrated that &gt;90% P.sub.i transport in cortical neurons, which displays similar kinetic parameters to those reported for cultured kidney proximal tubule epithelial cells and membrane vesicles, are sodium dependent and that this Na.sup.+ -dependent transport system is regulated through a Na.sup.+ -dependent P.sub.i cotransporner. B. Ni, et al., Proceedings of the National Academy of Sciences (USA), 91:5607-5611 (1994).\nThe present invention describes the cloning and characterization of a human brain Na.sup.+ -dependent P.sub.i cotransporter which is selectively expressed in discrete populations of neurons and glia. Fluorescent in situ hybridization (FISH) analysis demonstrates that this Na.sup.+ -dependent P.sub.i cotransporter is located in chromosome 19 (19q13.3) which has been linked to susceptible gene(s) for late onset Alzheimer's disease. M. Mullan and F. Crawford, Trends in Neurological Sciences, 16, 398-403 (1993). The characterization and treatment of physiological disorders is hereby furthered."}
{"text": "In phase contrast imaging, an object is radiated with coherent X-ray radiation, for example achieved by placing a grating structure in front of a conventional X-ray tube. For example, WO 2011/070521 A1 relates to differential phase contrast imaging and describes a respective system. A grating in front of the focal spot is provided to enhance the coherence length of the generated X-rays to a useful level. The grating is required to have a transparency that is reduced due to the requirement of having a small slits-to-pitch ratio, for example for the benefit of improved detection of phase shifts. However, it has been shown that for achieving sufficient image data quality, X-ray tubes with increased tube power are needed, which can result in expensive tubes. An example for increasing X-ray tube power is the provision of liquid metal jets acting as the anode for generating the X-ray radiation. For example, U.S. Pat. No. 6,995,382 B2 describes an arrangement for generating intensive radiation based on plasma generation, where the target generator has a multiple channel nozzle with a plurality of target jets for the generation of intensive short wave radiation. The plasma generated from the target jets merge into one extended plasma, leading to a powerful light source. However, the generation of plasma reduces the suitability for phase contrast X-ray imaging. It must be noted further that a source grating would have to be provided, which means additional effort due to the necessary manufacturing steps for the grating. Further, since the X-ray absorbing portions may be provided by gold material which is becoming a more and more expensive material, this also implies negative economic effects."}
{"text": "Mobility and a wideband has become a development trend of modern communication technologies, and the 3rd Generation Partnership Project (3GPP) has been devoted to a long Term Evolution (LTE) system evolving from a 3G system with the aim of developing the 3GPP radio access technologies toward a high data ratio, a short delay, and optimized packet data applications. The Multi-input Multi-Output (MIMO) technology at the physical layer has become one of crucial technologies in the current mobile communication system due to a number of advantages thereof, for example, the capacity of the system is extended through space division multiplexing of multiple antennas, the throughput of the system is improved using a multiplexing gain of multiple antennas, etc.\nAs the capacity is extended, interference also becomes higher accordingly due to the correlation between channels in a channel matrix, and if a base station can obtain somehow some Channel State Information (CSI) (which can be an instantaneous value, or can be short-term or middle- to long-term statistic information, then it may optimize power, rates, and even transmit directions applied to respective data streams in some preprocessing scheme, and possibly a part or all of interference between the data streams may be eliminated in advance at a terminal through preprocessing, thus achieving higher performance.\nAccordingly it is highly desirable in the industry to improve the extent of matching between the CSI fed back by the terminal, and a channel state."}
{"text": "Traditionally, footwear has been secured to the foot by means of a strip of material or thong. Early man discovered that removal and resecuring could be made relatively simple by forming holes near the edges of the material forming the footwear and threading the strip of material alternately through the holes of adjacent edges. Thus, was invented the first \"shoelace\" which has changed remarkably little in the intervening thousands of years.\nNot only does the shoelace provide the necessary function, loafers and the like aside, of ensuring satisfactory securing onto the foot, but it provides important decoration as well. Laces are made in a wide variety of patterns and colors to provide the wearer with an opportunity to use the lace as an accessory. And there is no doubt that these traditional laces perform these functions or they would not still be in near universal use after all these years.\nHowever, the traditional lace does have several drawbacks, notwithstanding its virtues and near universality. One of those is the limited area of the lace which is available for view and hence the limited visual impact a lace decoration can have. It is true that the lace width can be increased somewhat, but increasing the width makes drawing the lace through the holes more difficult. Also, the cross-section of the lace is invariably compressed and the lace deformed under tension, reducing its decorative impact.\nSecondly, threading the laces through the holes is a time consuming operation requiring manual dexterity. It is impractical to change the laces as often as a fashion conscious person might change other accessories, and yet given the option that is exactly what many, particularly young, people would very much like to do.\nThirdly, the fact that laces are so universal establishes a need for an alternative from a fashion standpoint. After all, it is the ability to choose which is fundamental to fashion."}
{"text": "Alkylation using HF is a widely used commercial refining and petrochemical process. Generally, alkylation is the addition of an alkyl group to another hydrocarbon. Commercially, HF is used to alkylate isobutane with propylene, butylene, and amylene isomers to produce high octane gasoline blending components, as well as to alkylate benzene in the manufacture of detergents.\nThe commercial process for alkylating isobutane with propylene is typical of many commercial HF alkylation process. Briefly, propylene and a stoichiometric excess of isobutane contact HF in a reactor where alkylation reactions produce alkylate. The reactor effluent passes to a settler in which the effluent separates into an acid phase containing HF and a hydrocarbon phase containing unreacted isobutane and alkylate. The acid phase is recycled to the reactor. The hydrocarbon phase passes to a fractionation column, commonly called an isostripper or a deisobutanizer, from which the isobutane is also recycled to the reactor and alkylate is recovered as product.\nIn addition to producing alkylate, HF alkylation processes also produce a by-product called acid-soluble oil, or ASO, which is soluble in HF. ASO is a recognized term in the art of alkylation, and is sometimes referred to as HF-soluble oil, polymer, conjunct polymer, or polymer by-product. Because ASO is soluble in HF, ASO in the settler concentrates not in the hydrocarbon phase but in the acid phase, and consequently the ASO is recycled to the reactor. In the reactor, ASO generally decreases the catalytic activity of the HF. While this effect can be beneficial at relatively low concentrations of ASO, if allowed to accumulate to relatively high concentrations ASO has a detrimental effect on the process. For this reason, ASO is typically removed at least periodically from a slip stream of the acid phase by a process that is commonly called regeneration.\nRegeneration is typically performed in two ways, internal regeneration and external regeneration. In internal regeneration, a slip stream of the acid phase, which is hereinafter referred to as a by-product stream, remixes with the hydrocarbon phase and passes to the isostripper. The ASO, being a heavy polymerized by-product, tends to pass downward in the isostripper and is rejected from the bottom of the isostripper and out of the process. The HF that enters the isostripper in internal regeneration tends to pass upward in the isostripper, is recovered from the overhead of the isostripper, and is recycled to the reactor. External regeneration, on the other hand, uses a stripping column commonly called an external regenerator to remove ASO from a by-product stream of the acid phase. In an external regenerator, either heat or isobutane or both strip HF overhead, leaving a residue of mostly ASO which is rejected from the bottom of the external regenerator and out of the process. The overhead of the external regenerator either is recycled as a condensed liquid to the reactor or is passed as a vapor to the isostripper where it may be used to assist in fractionation.\nA growing and recent trend in HF alkylation processes is the use of complexing agents to minimize the volatility of the acid phase. Because HF is volatile and the environmental risks arising from an accidental release of HF to the atmosphere are more and more undesirable, complexing agents that reduce the vapor pressure of HF and the tendency of HF to form an aerosol are now in greater use than ever before. Complexing agents typically contain at least one Lewis base site, which commonly comprises a Group 5A element such as nitrogen. Certain nitrogen-containing compounds, such as pyridine, picolines, quinoline, trimethylamine, and triethylamine, are known to form complexes with HF to reduce its volatility. For example, see Japanese Patent Disclosure No. 57(1982)-92502 (Oda et al). U.S. Pat. No. 5,073,674 (Olah) discloses that mixtures of HF and preferred nitrogen-containing compounds (complexing agents) such as ammonia, methylamines, ethylamines, propylamines, butylamines, pentylamines, pyridine, picolines, melamine, and hexamethylene-tetramine remained effective catalysts in alkylation of alkanes by alkenes. Mixtures of Olah's preferred complexing agents with HF are hereinafter referred to as HF-amine complexes.\nA problem that arises with the use of some complexing agents in HF alkylation processes is that both internal and external regeneration tend to reject not only ASO but also the HF-agent complex from the process. As a solution to this problem, a two-step method for removing ASO from a by-product stream containing HF, ASO, and HF-agent complex, where the complexing agent contains at least one Lewis base site containing a Group 5A element and the by-product stream has a molar ratio of HF per Lewis base site substantially above 5:1, is practiced in the art. First, a portion of the HF in the by-product stream is selectively removed to produce an HF-depleted stream having a molar ratio of HF per Lewis base site of 3:1 to 5:1. Then, the HF-depleted stream is separated into a hydrocarbon phase enriched in ASO and an acid phase depleted in ASO that contains a substantial portion of the HF-agent complex. The selectively removed HF and the acid phase are both recycled to the hydrocarbon alkylation step, and the hydrocarbon phase is rejected from the process. Thus, this method rejects ASO but not HF-agent complex from the process.\nOne of the problems with the method described in the previous paragraph is that the HF-enriched stream that contains the selectively removed HF also contains 10 to 15 wt-% ASO. The ASO in the HF-enriched stream tends to be lighter than the heavier ASO that remains in the HF-depleted stream and which is ultimately rejected from the process in the hydrocarbon phase. Nevertheless, like the heavy ASO, this light ASO which is recycled to the alkylation reactor has a detrimental effect on the alkylation reactions. Thus, a method is sought for rejecting ASO, in particular light ASO, from an HF alkylation process that uses an HF-agent complex, wherein the complexing agent contains at least one Lewis base site containing a Group 5A element.\nU.S. Pat. No. 5,073,674 (Olah) discloses a process of alkylating aliphatic hydrocarbons with alkenyl hydrocarbons in the presence of liquid HF-ammonia or HF-amine complexes.\nU.S. Pat. No. 5,191,150 (Child et al.) discloses a method of separating a stream containing ASO, HF, and sulfolane by first separating out HF and then gravitationally separating the HF-depleted stream into a polymer-rich stream and a sulfolane-rich stream.\nU.S. Pat. No. 5,382,746 (Child et al.) discloses an HF alkylation process comprising passing a reactor effluent settler acid phase containing sulfolane, ASO, and HF to an external regenerator; passing the external regenerator overhead stream containing HF and light ASO to an isostripper; recycling an isostripper overhead stream containing HF to the alkylation reactor; and recovering an isostrpper bottom stream containing alkylate and light ASO.\nIn motor fuel alkylation processes, the passage of the overhead vapor stream of an external regenerator to an isostripper, wherein the overhead vapor stream enters the isostripper at an intermediate point that is below the feed point of the isostripper, is practiced commercially.\nThe regeneration of HF used as catalyst is also practiced in the production of detergents by the reaction of C.sub.8 -plus normal olefinic hydrocarbons and benzene. The use in such a process of an external regenerator for regenerating HF and of another stripping column, which is called an HF stripper, for separating HF from the hydrocarbon-containing reactor effluent phase is shown in U.S. Pat. Nos. 4,237,327 (Winter, III) and 4,237,328 (Winter, III). U.S. Pat. No. 4,237,327 teaches a process in which an overhead vapor stream having a low concentration of ASO exits an external regenerator and enters an HF stripper at an upper intermediate point that is below the feed point of the HF stripper. U.S. Pat. No. 4,237,328 teaches passing the overhead vapor stream of an external regenerator and the overhead vapor stream of an HF stripper to a common single overhead condenser."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the examination of biological materials and, more particularly to providing a wavelength-tunable excitation light to excite a plurality of fluorophores associated with biological probe arrays and detect emissions from each of the fluorophores in response to the excitation light.\n2. Related Art\nSynthesized nucleic acid probe arrays, such as Affymetrix® GeneChip® probe arrays, and spotted probe arrays, have been used to generate unprecedented amounts of information about biological systems. For example, the GeneChip® Human Genome U133 Plus 2.0 probe array available from Affymetrix, Inc. of Santa Clara, Calif., is comprised of a single microarray containing over 1,000,000 unique oligonucleotide features covering more than 47,000 transcripts that represent more than 33,000 human genes. Analysis of expression data from such microarrays may lead to the development of new drugs and new diagnostic tools."}
{"text": "Dinettes are common features in recreational vehicles. A typical dinette includes a pair of bench seats and a table supported by one or more table posts or legs. Because space is limited in recreational vehicles, dinettes often can be converted into a bed by lifting the tabletop from its support posts, removing the support posts, and inserting the tabletop between the dinette seats on small ledges on the seats. The problem with conventional convertible dinettes is that the support posts or legs must be stowed when used as a bed."}
{"text": "1. Field of the Invention\nThis invention relates to an educational device, more particularly to a modular educational device that includes modular segments that can be assembled and disassembled to form a geometric shape.\n2. Description of the Related Art\nCrawl-through geometric units are used to help young children develop mobility abilities and identify geometric shapes. However, since such geometric units are generally made of wood, their shapes and sizes are fixed. As such, storage of the conventional geometric units is also inconvenient."}
{"text": "U.S. Pat. No. 4,625,712 and copending application Ser. No. 129,7l4 filed 12/07/87 disclose miniature high-speed blood pumps which can be threaded through a blood vessel to provide heart assist in emergency situations without major surgery. Typically, such pumps are inserted through, e.g., the femoral artery. An inflow cannula is positioned ahead of the pump during insertion, and this cannula must typically be pushed through the aortic valve in a retrograde direction.\nIn the past, such retrograde insertion was usually done by means of a wire guide over which the inflow cannula was slipped, and which was subsequently withdrawn. This method was unsatisfactory not only because it was awkward and carried a risk of injury to the vascular system and to the aortic valve, but also because it required the continuing observation of the wire during insertion by x-ray or other procedures. In an emergency, appropriate x-ray equipment may not be immediately available, and heart assist may have to be provided so quickly that the slow wire guide method may not be suitable.\nA requirement therefore exists for an inflow cannula which can be attached to the intake of an intravascular blood pump and can be pushed ahead of the pump for blind retrograde insertion into the left ventricle through the aortic valve without any additional apparatus and without danger of injury to the patient."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical scanning device for an image forming apparatus which scans an original document for forming a light image of the original document on an image carrier.\n2. Description of the Related Art\nA conventional optical scanning device for an image forming apparatus, such as a copying machine, is reciprocally moved by a driving mechanism to scan an original document. A light image is irradiated on an image carrier, for instance, a photoconductive drum, by the scanning of the original document with the optical scanning device.\nIn a first conventional optical scanning device, the device comprises a pair of shafts, a carriage on which a lamp, a reflector and a mirror are mounted, a pair of pipe-shaped supporting members mounted on the carriage, and a moving mechanism for moving the carriage. The pair of shafts is inserted into the pair of supporting members, and thus the pair of supporting members slidably moves along the shafts. The carriage is operatively connected to the moving mechanism to be reciprocally moved along the shafts via the pair of pipe-shaped supporting members.\nIn a second conventional optical scanning device, the device comprises a pair of guide rails, a carriage on which a lamp, a reflector and a mirror are mounted and having a pair of extended portions, a pair of slide members made of plastic with a low coefficient of friction mounted on the pair of extended portions, and a moving mechanism for moving the carriage. The carriage is movably placed on the pair of guide rails via the slide members. The carriage is operatively connected to the moving mechanism to be reciprocally moved along the rails via the pair of slide members.\nHowever, the above-mentioned conventional optical scanning devices have many problems. That is, force moments occur about the leading edge in the moving direction of the carriages at the start of forward movement and at the start of backward movement during the scanning operation of the optical scanning devices. In the first conventional optical scanning device, although a certain amount of clearance is required between the shaft and the pipe-shaped supporting member, when a force moment occurs at the start of reciprocal movement, rattling occurs between the pair of shafts and the pipe-shaped supporting members due to this clearance. This leads to a problem of poor picture quality due to the occurrence of image blurring.\nIn the second conventional optical scanning device, since the carriage is held only on the guide rails by its own weight via the slide members, there is a problem in that the carriage is liable to jump up at the start of reciprocal movement or due to other external forces and thus image blurring will occur. In either case, the higher the copying speed, the greater the blurring, and thus there is a problem in that the system cannot be applied to high-speed models."}
{"text": "The present invention relates to fertilization wherein irrigation waters of agricultural irrigation systems contain fertilizer nutrients.\nThe agriculture industry adds fertilizers to the plant environs, such as the soil, to enhance crop growth and subsequent yields. Commercial fertilizers are usually selected of a variety of formulations depending on the crop and its nutrient requirements.\nFertilizers generally are classified according to their NPK content. NPK is common terminology used in the fertilizer industry and stands for: (1) N—the amount of nitrogen in the formulation as N; (2) P—the amount of phosphorus in the formulation as P2O5; and (3) K—the amount of potassium in the formulation as K2O. Nitrogen, phosphorus and potassium are the basic plant nutrients or macronutrients that are taken up and utilized by the growing crops, and they are commonly provided or augmented by the addition of fertilizers (NPK fertilizers).\nA fertilizer, as that term is used herein and as generally understood, refers to the nutrient-containing materials used to deliver fertilizer nutrients to a crop. Conventional fertilizers typically will contain materials that are extraneous to the crop's nutrient-uptake and soil condition (“yield-extraneous constituents”) but which, for practical and/or other reasons, are necessary to the delivery of the nutrients. The process of delivering fertilizer nutrients to crops is referred to as fertilization.\nThe fertilization method of adding fertilizers to the water being used to irrigate the crops is called “fertigation”, reflecting this combination of irrigation and fertilization. Fertigation reduces equipment, fuel and labor expended in the addition of fertilizers in comparison to mechanical delivery of fertilizers to the crop, and thus fertigation achieves a significant overall cost savings.\nTo conserve water, which is decreasing in availability and increasing in cost, current conventional technology includes micro-irrigation systems that deliver precise amounts of water directly to the soil holding the root system of the plant that is being grown. In the past twenty to thirty years, a large percentage of crop producers in the western and southwestern United States have converted from flood and sprinkler irrigation systems to micro-irrigation technology. Micro-irrigation systems contain devices called emitters, micro-sprinklers or other such devices that provide the precise amounts of water directly to the desired soil site, namely the soil holding the roots of the plant or crop being irrigated.\nMicro-irrigation systems are sensitive to water quality and the inclusion of fertilizers and other additives stemming from the refinement of their micro components. These emitters, micro-sprinklers or other micro devices have very tiny orifices and/or a long tortuous narrow passageway that provide the requisite pressure for delivery of precise amounts of water in a uniform manner to each plant in the crop being irrigated so long as deposits do not build up inside them. Deposits from any source foul or plug these micro components. The smallest particle or foreign material can cause fouling of these devices. Water quality and the inclusion of fertilizers can, and frequently does, cause severe plugging problems. The problems arise from a number of factors: (1) the irrigation water is typically obtained from wells, reservoirs, canals, lakes, or rivers which contain various amounts of dissolved minerals; and (2) fertilizers can form insoluble salts and/or cause particulate formation when added to the water. Macro-irrigation systems (macro-sprinkler irrigation systems) mainly tolerate these conditions, while micro-irrigation systems are extremely intolerant.\nIn more detail, the addition of conventional fertilizer formulations to irrigation water normally increases the loading of inorganic salts over that already in the water. When the loading, or the combined loading, is too high, the solubilities of at least some of the naturally-occurring irrigation-water minerals and/or added compounds are exceeded and particulate formation increases dramatically. When particulates form, significant deposits begin to build up throughout the irrigation system. The end result for emitters or micro-sprinklers is plugging.\nPlugging results in uneven distribution of water and nutrients to the crop being irrigated, and in some cases, the complete shut-down of the micro-irrigation system. Problem-free use of additives such as fertilizers in micro-irrigation systems is normally seen only in irrigation systems that use relatively pure water sources.\nThe conventional agricultural practice is to make intermittent or periodic applications of fertilizers. Such intermittent or discontinuous additions might be a single addition, or a plurality of additions, of large amounts (high concentrations) of fertilizer during a brief time interval each growing season or crop cycle. (The number of applications per growing season or crop cycle usually depends on the crop and/or the type of fertilizer being added.) When the fertilizer-delivery method is fertigation, fertilizers are typically slug fed into the irrigation system as quickly as possible to minimize the labor requirements and ease material handling. Slug feeding of a block (portion of a field) normally entails feeding the large amounts (high concentrations) of fertilizer to the irrigation water over a six to seven hour period during irrigation, and then, after the fertilizer feed is shut off, continuing the irrigation of that block for an additional two to three hours to rinse out all of the fertilizer that is contained inside the irrigation system, insuring that all of the fertilizer intended for the block is in fact delivered to the block. When an entire field is irrigated on a single day, slug feeding does not require a post-fertigation rinsing period.\nContinuous fertigation methods, in which the irrigation water applied throughout a crop cycle, or during a prolonged period of a crop cycle, contains fertilizers, are known from U.S. Pat. No. 8,721,758, issued May 13, 2014, inventors Miller et al., the content of which is hereby incorporated by reference hereinto, and U.S. Pat. No. 8,986,417, issued Mar. 24, 2015, inventors Miller et al., the content of which is hereby incorporated by reference hereinto. Slug feeding fertigation methods are known from U.S. application Ser. No. 14/564,594, filed Dec. 9, 2014, inventors Miller et al., the content of which is hereby incorporated by reference hereinto.\nThe cost of commercial fertilizer formulations is, however, itself significant, and commercially viable fertilizer formulations (formulations sufficiently inexpensive for bulk agricultural use) typically include, as mentioned above, yield-extraneous constituents which do not contribute to plant nutrition or soil condition, and can even be undesirable components. Further, the bulk weight of commercial fertilizer formulations typically is water which increases the shipping costs. The water of a fertilizer formulation is analogous to, or at least approaches being, a yield-extraneous constituent because its contribution to the water requirements of a crop normally is negligible, particularly in comparison to its adverse impact on shipping, storage and handling but since water is normally required for an unimpeded container-to-irrigation-system fertilizer addition it is analogous to, but still not considered herein, a yield-extraneous constituent.\nThe terms “micro-irrigation” and “macro-irrigation” as used herein refer respectively to (1) micro-sprinklers, drip, and subsurface drip systems and (2) sprinkler systems without micro components which are primarily overhead sprinklers. The terms “overhead sprinkler” and “overhead sprinklers” as used herein refer to irrigation systems in which the irrigation water is emitted or sprayed from sprinkler heads, nozzles or other irrigation devices disposed at a position elevated from ground level, that normally (but not necessarily) are engaged directly above the crop being irrigated. Overhead sprinkler irrigation systems are normally macro-irrigation systems because micro-irrigation normally depend on a close proximity between the emitters and soil immediate the crop while overhead sprinklers are normally spaced apart from the soil immediate the crop. Such micro- and macro-irrigation systems are collectively referred to herein as “emitter-irrigation” systems because they each emit or eject sequential small quantities, or streams, of irrigation water from irrigation lines or tubes directed more or less towards the soil of crop being irrigated."}
{"text": "1. Field of the Invention\nThe present invention relates to vapor generators and, more particularly, to a vapor generator fluid flow control system, permitting controlled fuel, air and water flow and pressure regulation.\n2. History of the Prior Art\nVapor generators of the kind in which a fuel-air mixture is combusted in the direct presence of feed water to produce a useful mixture of steam and non-condensibles are known. Such vapor generators are shown in U.S. Pat. No. 4,211,071 issued to the assignee of the present invention. In accordance with that invention, a vapor generator is provided in which several interrelated means are employed to improve the quality of combustion in the generator so that a product steam substantially free of carbon monoxide results. In effecting this result, means are provided for dividing the air feed into two parts which are both delivered to the generator for combustion. Such methods and apparatus are effective in producing a well-mixed stoichiometric mixture of fuel and air to provide completeness of combustion and reducing production of carbon monoxide to extremely low levels. Such an accomplishment is a marked advance over the prior art and illustrates the emphasis placed on controlled flow rates and mixture conditions.\nThe prior art is also replete with apparatus manifesting high temperature product vapor having high carbon monoxide content. Such conditions are objectionable for many applications and dangerous for some of them. High carbon monoxide production is traceable to incomplete combustion. This is, in turn, traceable in part to difficulties in maintaining a stable lean flame and, in part, to excessive quenching of the flame through direct radiation and convective contact between the flame and the feed water. Once a stoichiometric mixture is obtained, it is, therefore, necessary to maintain the combustion level without varying any of the three major constituents of fuel, air, or water. Variations would affect the combustion and the carbon monoxide production therefrom.\nOther problems occur in certain ones of the prior art vapor generators when they are operated at low pressures. Conventionally, as shown in U.S. Pat. No. 4,288,978, a vapor generator may be especially adapted for low pressure operation by the provision of certain features such as a pilot burner for striking a stable flame. It may be seen that low pressure operation, likewise, requires a specific reduction in the air, fuel and water input for a low volume output from the vapor generator. Disproportionate variances in any of the three constituents will seriously affect the combustion level and the resulting heat and vapor content of the downstream product.\nFew prior art systems effectively address dual output vapor generator operation. The ability to fine tune the fuel, air and water flow for a single vapor generator output is critical enough. Once this critical equilibrium is reached, it is tantamount to starting over to vary the generator operation to another flow. Varying the generator volume necessitates adjustment in the air, fuel and water flow rates which will directly affect both the heat and the water vapor level of the combustion products. When such generators are used in industrial applications necessitating precise control of heat and vapor, these variations are not tolerable. One such application is the curing of concrete where the temperature and moisture level of the vapor products have a direct effect on the curing process and resulting structural condition of the cast product.\nThe use of steam for concrete curing is not novel in and of itself. Steam has been produced from conventional boilers for such methods and processes, but the cost of a boiler operation is much higher than that of a vapor generator type device. Moreover, the temperature and vapor level of the steam product can be more closely controlled than with a conventional boiler. Moreover, super heated vapor can be provided for specific types of aggregate to be cured. The flexibility of the vapor generator also permits substantially instantaneous operation as compared to the threshold period for a boiler method. The obvious disadvantage is the substantially single flow rate afforded by most conventional vapor generators. Since any number of kilns may be run at one time for particular production schedules, excess vapor from a single flow rate vapor generator system would, thereby, be wasted.\nConventional techniques for varying the output of standard vapor generators includes reducing the rate of combustion and water flow while maintaining a constant air flow rate. The advantage of such a system is simplicity and cost in that reliable variations in air flow rates have been, to date, difficult to attain. The paramount disadvantage to the continuous air flow volume is the excess air ingressing into the kiln and the adverse effects to the curing process which affects the ultimate humidity level within the kiln and imparts non-uniform cooling characteristics. Variations in the air flow volume have, to date, been addressed by multiple speed blowers and throttle systems which either decrease the air intake to the blower or exhaust therefrom. Unfortunately, the relatively large motors and blower units necessary for vapor generator volumes generate large amounts of heat which the air volume dissipates. When the air volume is throttled for constant motor speed, heat dissipation is inhibited and variations in flow rate result as well as having degenerative effects on the motor and the blower unit. The obvious alternative to such flow problems is a multiple speed motor, but the cost and availability for such systems are generally disproportionate to the vapor generator construction.\nDual output vapor generator applications also require dual flow feedwater systems which may be accurately set in preselected flow configurations as well as precisely maintained. Precise fluid flow is an integral element of the aforesaid combustion efficiency particularly in stoichiometric mixtures. The primary area of fluid flow consideration in such vapor generator systems is in the water flow controlling network. Prior art flow control devices though capable of preselect flow volumes are generally not sensitive to variations in flow pressure which affects the flow volume.\nIt would be an advantage therefore to overcome the problems of the prior art by providing a variable output vapor generator having a constant volume fluid flow system which affords automatic operation between high and low fluid flow rates and which is sensitive to fluctuations in fluid flow pressure. One approach to dual output vapor generators and a discussion of prior art problems associated therewith is discussed in co-pending U.S. patent application Ser. No. 554,780 assigned to the assignee of the present invention. The method and apparatus of the present invention comprises an advanced system overcoming the problems set forth above by providing a combination of a pressure responsive valve and fine adjustment valve in parallel flow communication. This parallel flow system is itself aligned in parallel flow communication with a second matching system adapted for selective actuation for high and low volume operation. In this manner both the high volume and low volume flow rates within the water flow system of a vapor generator may be pre-adjusted for automatic actuation as needed by demand conditions. This has been done in a configuration facilitating automatic temperature control and which is pressure sensitive for accurate volumetric control without adversely affecting the efficiency of the generator system."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and device for measuring interface trap density in a semiconductor device.\n2. Description of the Related Art\nAs a semiconductor technology develops constantly, the length of gates gradually become shorter and the thickness of gate oxide layers gradually become thinner. Accordingly, characteristics of semiconductor devices are inevitably affected by these smaller dimensions. One of the characteristics affected is the interface trap density between a silicon substrate and a gate oxide layer.\nGenerally, when a gate oxide layer is grown, a sufficient bonding is not made between a silicon atom and an oxygen atom, so that a dangling bond (where there are insufficient oxygen atoms) is generated at an interface between a silicon substrate and the gate oxide layer.\nA dangling bond easily attracts an electron when a transistor operates, thereby increasing the interface trap density of the gate oxide layer. Accordingly, the quality of the gate oxide layer deteriorates or the driving current is reduced, which may degrade some characteristics of a semiconductor device.\nTraditionally, a charge pumping test has been used as a method for measuring a surface state located under a gate oxide layer. The interface trap density may be calculated from data obtained using this test. However, when the conventional method is applied to a device having a very thin gate oxide layer, gate tunneling leakage current and/or quantum mechanical effects can result in an incorrect calculation of the interface trap density (generally resulting in a calculated interface trap density that is too high)."}
{"text": "Control moment gyroscopes (CMGs) are commonly employed in satellite attitude control systems. A generalized CMG may include a main CMG housing and an inner gimbal assembly (IGA), which is mounted within the main CMG housing for rotation about a gimbal axis. The IGA includes, amongst other components, a rotor assembly, at least one spin motor, and a rotor support structure. In one known design commonly referred to as a “shell rotor architecture” or, more simply, a “shell rotor,” the rotor assembly includes an annular rotor rim disposed between first and second hemi-spherical rotor shell pieces. The rotor rim is welded to the first and second rotor shell pieces along first and second circumferential rim-shell weld joints, respectively. Collectively, the rotor rim, the first rotor shell piece, and the second rotor shell piece form a generally spherical body having a central channel through which a rotor shaft extends. During operation, the spin motor rotates the rotor assembly about a spin axis at a relatively high rate of speed. The IGA is, in turn, selectively rotated about the gimbal axis, which is generally perpendicular to the rotor spin axis, by an actuator (e.g., a torque module assembly) mounted within the main CMG housing. When the IGA is “gimbaled” in this manner, a controlled torque is generated about an output axis normal to the spin and gimbal axes due to momentum transfer of the rotor assembly. This gyroscopic torque is imparted to the host satellite through a CMG mount interface, such as an annular flange bolted to the exterior of the satellite. By selectively rotating the IGA about the gimbal axis, the satellite's attitude may be adjusted and re-adjusted, as needed, in a highly controlled manner.\nCMG shell rotors of the type described above have been developed and commercially implemented with considerable success by the assignee of the present Application, Honeywell International Inc. In general, CMG shell rotors provide superior performance characteristics as compared to other conventionally-known CMG rotors, including superior momentum-per-weight ratios and relatively high operational speed limits, as well as high angular stiffness. However, further improvements in momentum-per-weight ratios, speed capabilities, and other measures of CMG performance (e.g., operational lifespan) are still desirable. Moreover, as a notable disadvantage, the manufacturing schedule and costs associated with manufacture of CMGs having shell rotor-based architectures tend to be substantial. It is not uncommon for commercial production of a shell rotor to cost upwards of several hundred thousand U.S. dollars and to require over one year to complete.\nIt is thus desirable to provide CMGs including advanced rotors having improved performance characteristics (e.g., increased momentum-per-weight ratios, higher speed limits, longer operational lifespans, and so on) as compared to shell rotor assemblies and other conventionally-known CMG rotor assemblies. Ideally, embodiments of such a CMG, and specifically embodiments of the advanced rotor included within the CMG, would be amenable to production utilizing cost effective and time efficient manufacturing processes. It would also be desirable to provide manufacturing processes for producing CMGs including such high performance rotors. Other desirable features and characteristics of the present invention will become apparent from the subsequent Detailed Description and the appended claims, taken in conjunction with the accompanying drawings and this Background."}
{"text": "1. Field of the Invention\nThe present invention relates to metering or reporting of energy consumption, and, more particularly, to enabling temporary associations of parent energy consumption meters and submeters.\n2. Related Art\nIn a number of circumstances, it is useful for an energy purchaser to use a fixed arrangement of submeters to measure selected portions of the purchaser's overall energy consumption. (Meters for measuring energy consumption are well known. See, for example, http://www.net-metering.com/power_meters.html, http://www.watthourmeters.com, and, in particular, http://www.watthourmeters.com/modern/centron.html.) Measuring electrical use in apartment complexes is one such submetering application that is common. It is economical to use fixed electrical submeters in this situation because each apartment's consumption is large and the arrangement of submeters remains fixed for a long time. Submetering is generally not economical in situations where the total amount of energy consumption is small, the portion of overall energy consumption that is of interest is small or the submetering would be temporary. For example, although patrons commonly plug notebook computers and the like into electrical outlets in public places like libraries, book stores, coffee shops, hotel lobbies and airports, these outlets are typically not submetered because they supply only a little energy and also because patrons use them for only a short while.\nA new situation is developing in which energy consumption is moderate, widespread and highly temporary. That is, the number of hybrid automobiles is increasing dramatically. These automobiles may need to be electrically recharged in public places, such as at parking garages or street side parking spaces. While this could be facilitated by installing electrical outlets having coin-operated timers, this tends to be inconvenient for the user. Thus, a need exists for improvements in energy consumption submetering."}
{"text": "1. Field of the Invention\nThe present invention relates generally to magnetoresistive (MR) sensor elements employed within magnetic data storage and retrieval. More particularly, the present invention relates to enhanced magnetoresistive (MR) resistivity sensitivity anisotropic magnetoresistive (MR) sensor elements employed within magnetic data storage and retrieval.\n2. Description of the Related Art\nThe recent and continuing advances in computer and information technology have been made possible not only by the correlating advances in the functionality, reliability and speed of semiconductor integrated circuits, but also by the correlating advances in the storage density and reliability of direct access storage devices (DASDs) employed in digitally encoded magnetic data storage and retrieval.\nStorage density of direct access storage devices (DASDs) is typically determined as areal storage density of a magnetic data storage medium formed upon a rotating magnetic data storage disk within a direct access storage device (DASD) magnetic data storage enclosure. The areal storage density of the magnetic data storage medium is defined largely by the track width, the track spacing and the linear magnetic domain density within the magnetic data storage medium. The track width, the track spacing and the linear magnetic domain density within the magnetic data storage medium are in turn determined by several principal factors, including but not limited to: (1) the magnetic read-write characteristics of a magnetic read-write head employed in reading and writing digitally encoded magnetic data from and into the magnetic data storage medium; (2) the magnetic domain characteristics of the magnetic data storage medium; and (3) the separation distance of the magnetic read-write head from the magnetic data storage medium.\nWith regard to the magnetic read-write characteristics of magnetic read-write heads employed in reading and writing digitally encoded magnetic data from and into a magnetic data storage medium, it is known in the art of magnetic read-write head fabrication that magnetoresistive (MR) sensor elements employed within magnetoresistive (MR) read-write heads are generally superior to other types of magnetic sensor elements when employed in retrieving digitally encoded magnetic data from a magnetic data storage medium. In that regard, magnetoresistive (MR) sensor elements are generally regarded as superior since magnetoresistive (MR) sensor elements are known in the art to provide high output digital read signal amplitudes, with good linear resolution, independent of the relative velocity of a magnetic data storage medium with respect to a magnetoresistive (MR) read-write head having the magnetoresistive (MR) sensor element incorporated therein.\nWithin the general category of magnetoresistive (MR) sensor elements, anisotropic magnetoresistive (MR) sensor elements are presently of considerable interest insofar as anisotropic magnetoresistive (MR) sensor elements have found widespread use in various magnetic data storage and transduction applications, and in particular digitally encoded magnetic data storage and transduction applications. Within the context of the present disclosure, anisotropic magnetoresistive (MR) sensor elements are intended as magnetoresistive (MR) sensor elements fabricated employing either a single magnetoresistive (MR) sensing layer or a multiplicity of magnetoresistive (MR) sensing layers which are separated by dielectric spacer layers, and where a magnetoresistive (MR) effect within such an anisotropic magnetoresistive (MR) sensor element derives from magnetically encoded data flux perturbation of a magnetically aligned domain within either the single magnetoresistive (MR) sensing layer or the multiplicity of magnetoresistive (MR) sensing layers. Such anisotropic magnetoresistive (MR) sensor elements include, but are not limited to single stripe magnetoresistive (SSMR) sensor elements, soft adjacent layer (SAL) magnetoresistive (MR) sensor elements and dual stripe magnetoresistive (DSMR) sensor elements.\nWithin magnetoresistive (MR) sensor element fabrication, including but not limited to anisotropic magnetoresistive (MR) sensor element fabrication, it is particularly desirable to fabricate magnetoresistive (MR) sensor elements with enhanced levels of magnetoresistive (MR) resistivity sensitivity. For purposes of clarity, and within the context of the present disclosure, magnetoresistive (MR) resistivity sensitivity is intended as a measure of proportion of resistance change normalized to an absolute resistance of a magnetoresistive (MR) sensor element (i.e. dR/R) when measuring a magnetic signal within a magnetic data storage medium while employing the magnetoresistive (MR) sensor element. The magnetoresistive (MR) resistivity sensitivity of a magnetoresistive (MR) sensor element is alternatively known as the magnetoresistive (MR) coefficient of the magnetoresistive (MR) sensor element. Magnetoresistive (MR) sensor elements exhibiting enhanced magnetoresistive (MR) resistivity sensitivity are desirable within the art of magnetoresistive (MR) sensor element fabrication since such enhanced magnetoresistive (MR) resistivity sensitivity clearly inherently allows for detection within a magnetic data storage medium of weaker magnetic signals with increased linear density and thus also inherently allows for an increased areal density of the magnetic data storage medium within a magnetic data storage enclosure which employs the magnetoresistive (MR) sensor element which exhibits the enhanced magnetoresistive (MR) resistivity sensitivity.\nIt is thus towards the goal of forming for use within magnetic data storage and retrieval anisotropic magnetoresistive (MR) sensor elements, such as but not limited to single stripe magnetoresistive (SSMR) sensor elements, soft adjacent layer (SAL) magnetoresistive (MR) sensor elements and dual stripe magnetoresistive (DSMR) sensor elements, with enhanced magnetoresistive (MR) resistivity sensitivity, that the present invention is directed.\nVarious magnetic sensor elements, including but not limited to anisotropic magnetoresistive (MR) sensor elements, which possess desirable properties have been disclosed within the art of magnetic sensor element fabrication, including but not limited to anisotropic magnetoresistive (MR) sensor element fabrication.\nFor example, Goubau et al., in U.S. Pat. No. 5,268,806, disclose a magnetoresistive (MR) sensor element having a conductor lead structure which remains stable not only during processing when fabricating the magnetoresistive (MR) sensor element, but also over the useful operational life of the magnetoresistive (MR) sensor element. The magnetoresistive (MR) sensor element employs a conductor lead structure comprising: (1) a seed layer contacting a magnetoresistive (MR) layer within the magnetoresistive (MR) sensor element, where the seed layer is formed from a material selected from the group consisting of chromium, tungsten, an alloy of titanium and tungsten, and an alloy of tantalum and tungsten, and where the seed layer has a body centered cubic lattice structure; and (2) a conductor layer formed upon the seed layer, where the conductor layer is formed of tantalum, and where the conductor layer also has a body centered cubic lattice structure.\nIn addition, Fontana Jr. et al., in U.S. Pat. No. 5,701,223, disclose a spin valve magnetoresistive (SVMR) sensor element with improved magnetoresistive (MR) properties, such as an enhanced magnetoresistive (MR) resistivity sensitivity, as well as improved chemical properties, such as an enhanced corrosion resistance of an antiferromagnetic pinning material layer within the spin valve magnetoresistive (SVMR) sensor element. The spin valve magnetoresistive (SVMR) sensor element realizes the foregoing objects by employing when forming the spin valve magnetoresistive (SVMR) sensor element: (1) a pair of anti-parallel antiferromagnetically biased ferromagnetic material layers separated by a non-magnetic conductor spacer layer to provide a composite pinned ferromagnetic layer within the spin valve magnetoresistive (SVMR) sensor element, in conjunction with; (2) a pinning material layer formed of a pinning material selected from the group consisting of nickel oxide, nickel-cobalt alloy oxides and iron-manganese-chromium alloys.\nFurther, Lee et al., in U.S. Pat. No. 5,731,936, disclose an anisotropic magnetoresistive (MR) sensor element, such as a soft adjacent layer (SAL) magnetoresistive (MR) sensor element, with an enhanced magnetoresistive (MR) resistivity sensitivity as well as an enhanced thermal stability. The anisotropic magnetoresistive (MR) sensor element employs a nickel-iron-chromium alloy or a nickel-chromium alloy spacer layer adjoining a nickel-iron alloy magnetoresistive (MR) layer formed within the anisotropic magnetoresistive (MR) sensor element.\nStill further, Chang et al., in U.S. Pat. No. 5,763,108, disclose a high saturation magnetization magnetic material for use when fabricating a magnetic head, and a magnetic head fabricated employing the high saturation magnetization magnetic material for use when fabricating the magnetic head. The high saturation magnetization magnetic material has an elemental composition comprising about 40 to about 60 weight percent iron, about 40 to about 60 weight percent nickel and about 0.002 to about 1 percent tin.\nFinally, Lin, in U.S. Pat. No. 5,768,071, discloses a spin valve magnetoresistive (SVMR) sensor element having an improved magnetic stability of a pinned ferromagnetic layer within the spin valve magnetoresistive (SVMR) sensor element. The spin valve magnetoresistive (SVMR) sensor element realizes the foregoing object by employing when forming the spin valve magnetoresistive (SVMR) sensor element a discontinuous nonmagnetic interlayer interposed between an antiferromagnetic pinning material layer and the pinned ferromagnetic layer within the spin valve magnetoresistive (SVMR) sensor element.\nDesirable in the art of magnetoresistive (MR) sensor element fabrication are additional methods and materials which may be employed for forming with enhanced magnetoresistive (MR) resistivity sensitivity anisotropic magnetoresistive (MR) sensor elements, including in particular dual stripe magnetoresistive (DSMR) sensor elements.\nIt is towards the foregoing object that the present invention is directed.\nA first object of the present invention is to provide a method for forming an anisotropic magnetoresistive (MR) sensor element, along with the anisotropic magnetoresistive (MR) sensor element formed employing the method.\nA second object of the present invention is to provide a method and an anisotropic magnetoresistive (MR) sensor element in accord with the first object of the present invention, where the anisotropic magnetoresistive (MR) sensor element has an enhanced magnetoresistive (MR) resistivity sensitivity.\nA third object of the present invention is to provide a method and an anisotropic magnetoresistive (MR) sensor element in accord with the first object of the present invention and the second object of the present invention, which method is readily commercially implemented.\nIn accord with the objects of the present invention, there is provided by the present invention a method for forming an anisotropic magnetoresistive (MR) sensor element, as well as the anisotropic magnetoresistive (MR) sensor element formed in accord with the method. To practice the method of the present invention, there is first provided a substrate. There is then formed over the substrate a seed layer, where the seed layer is formed of a magnetoresistive (MR) resistivity sensitivity enhancing material selected from the group consisting or nickel-chromium alloys and nickel-iron-chromium alloys. There is then formed over the seed layer a nickel oxide material layer. There is then formed over the nickel oxide material layer a magnetoresistive (MR) ferromagnetic layer.\nThe method of the present invention contemplates an anisotropic magnetoresistive (MR) sensor element fabricated in accord with the method of the present invention.\nThe present invention provides a method for forming an anisotropic magnetoresistive (MR) sensor element, and the anisotropic magnetoresistive (MR) sensor element formed employing the method, where the anisotropic magnetoresistive (MR) sensor element has an enhanced magnetoresistive (MR) resistivity sensitivity. The method of the present invention realizes the foregoing object by employing when forming the anisotropic magnetoresistive (MR) sensor element: (1) a seed layer formed of a magnetoresistive (MR) resistivity sensitivity enhancing material selected from the group consisting of nickel-chromium alloys and nickel-iron-chromium alloys, in turn having formed thereover; (2) a nickel oxide material layer, in turn having formed thereover; (3) a magnetoresistive (MR) ferromagnetic layer. While it is not entirely clear why an anisotropic magnetoresistive (MR) sensor element formed in accord with the present invention exhibits an enhanced magnetoresistive (MR) resistivity sensitivity, it is nonetheless clear that an anisotropic magnetoresistive (MR) sensor element formed in accord with the present invention exhibits an enhanced magnetoresistive (MR) resistive sensitivity in comparison with an analogous anisotropic magnetoresistive (MR) sensor formed otherwise in accord with the present invention, but absent a nickel oxide material layer formed interposed between: (1) a seed layer formed of a magnetoresistive (MR) resistivity sensitivity enhancing material selected from the group consisting of nickel-chromium alloys and nickel-iron-chromium alloys; and (2) a magnetoresistive (MR) ferromagnetic layer.\nThe method of the present invention is readily commercially implemented. The present invention employs methods and materials as are generally known in the art of magnetoresistive (MR) sensor element fabrication but employed within the context of specific locations and limitations within an anisotropic magnetoresistive (MR) sensor element to provide the present invention. Since it is a novel ordering and use of methods and materials which at least in part provides the present invention, rather than the existence of methods and materials which provides the present invention, the method of the present invention is readily commercially implemented."}
{"text": "FIG. 1 is a view illustrating a conventional ice tray assembly for a refrigerator. The conventional ice tray assembly for the refrigerator includes a casing 10 installed in the refrigerator, and an ice tray 20 placed in the casing 10, for making ices. A shelf 12 on which the ice tray 20 is put is formed in the casing 10. A user fills the ice tray 20 with water and puts the ice tray 20 into the casing 10, so that the ices are made in the ice tray 20 by cool air flowing in the refrigerator.\nHowever, in the conventional ice tray assembly for the refrigerator, when the casing 10 with the ice tray 20 mounted therein is installed in a freezing chamber for freezing food, the casing 10 occupies a large area of a freezing space in the freezing chamber. In the case that the casing 10 is installed at a door for opening and closing the freezing chamber, when the door is closed, water filled in the ice tray 20 may overflow out of the casing 10.\nMoreover, as a valve is driven in a state where a water tank is not completely mounted in the freezing chamber, some of water stored in the water tank may not be supplied to the ice tray 20 but dispersed into the freezing chamber."}
{"text": "In the automobile industry, vehicle cross-support members, crossbeams, cross car beams, and other similar components (collectively referred hereinafter as “vehicle cross-support members” or “cross-support members”) are utilized as part of the vehicular body structure. The vehicle cross-support member normally spans between or fastened to a pair of laterally disposed vertical pillars, or A-pillars, in the region generally below the windscreen and a cowl top, and between a forward engine compartment and a rearward passenger compartment, so as to extend in a direction transverse to the length of an automobile. As part of a motor vehicle body, the cross-support member provides for cross car stiffness and rigidity against for example side load impacts.\nLocated forwardly of the driver and the front-row passenger, the cross-support member also supports or provides mounting surfaces for various vehicle components, including an instrument panel, a glove and/or storage compartment, a center console, a dashboard and a steering column member. Although direct attachment may be possible, various attachment devices or assemblies could be utilized to effect attachment of such vehicle components to the cross-support member. The attachment devices and assemblies are preferably designed such that the noise and vibrations transferring from the vehicle engine, drivetrain, transmission or any other moving parts are reduced for optimal passenger comfort. For vehicle components subject to or exerting a load or force, the attachment devices and assemblies preferably incorporate constructions adapted for improving load transfer, dispersion and/or distribution. Reduction of weight and production costs are also of important consideration in designing and implementing the attachment devices and assemblies.\nA number of attachment or mounting assemblies for securing a steering column member in a vehicle are known. By way of an example, U.S. Pat. No. 8,256,830 to Hitz describes a plurality of standardized support arms, a number of which are configured in the form of steering column support arms each having a U-shaped cross section defined by a pair of legs. Each of the legs forming part of the steering column support arm of Hitz include an area secured in a preformed transversal slot of a vehicle cross-support member on the driver side. Hitz furthermore describes reinforcement members having the same U-shaped cross section, and which are coupled to the steering column support arms to form a rectangular cross section. Hitz suggests that the standardized support arms may alternatively be formed as other support or reinforcement components secured to the cross-support member, such as a transmission tunnel support arm and a splash shield support arm, so as to permit reduction of production costs associated with fabricating multiple vehicle components.\nU.S. Pat. No. 8,146,986 to Bierkamp describes a holder configured for attachment to a tubular vehicle cross-support member to effect securement of different vehicle components. The holder of Bierkamp possesses a U shaped configuration defined by two limbs interconnected by a bridge, and which are sized to be received in a parallel spaced slots transversely oriented to a longitudinal axis of the cross-support member. Bierkamp suggests that two or more holders may be secured to the cross-support member to effect steering column attachment. The specific configuration of the holder of Bierkamp is described as advantageous in that it allows for sliding and rotational adjustment or positioning of the holder on the tubular cross-support member prior securing or welding the holder thereto.\nThe attachment devices described in Hitz and Bierkamp may suffer disadvantages in that forces generated by a steering column member, such as rotational forces, may not be transferred or distributed efficiently to other structural elements of the vehicle. Furthermore, as the generated forces are mainly channeled to the vehicle cross-support member, deficiencies in the attachment between the cross-support member and the devices of Hitz and Bierkamp may possibly lead to structural and/or functional failures. Given the increasing complexity and components included with steering column members, such as airbags and various control and adjustment elements, improvements to the load transfer and distribution from the steering column member may be desired."}
{"text": "Ethanol fuel derived from plants is expected to be liquid fuel alternative to gasoline to prevent increase in carbon dioxide gas. When both sugar and ethanol are produced from a sugar juice derived from a plant, the following method has been employed. Specifically, first, sugar is produced from a sugar juice. The sugar juice after the sugar production is fermented by using a microorganism to produce ethanol (see, for example, JP-A 2004-321174)."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor fabrication techniques and, more particularly, to a fabrication technique for improving conductivity at an interface between a polysilicon layer and a conductive layer in an integrated circuit device.\n2. Description of the Related Art\nThis section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present invention, which are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present invention. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nSemiconductor memory devices, such as dynamic random access memory (DRAM) devices, are widely used for storing data in systems such as computer systems. A DRAM memory cell typically includes an access device such as a field effect transistor (FET) coupled to a storage device such as a capacitor. The access device allows the transfer of charged electrons to and from the storage capacitor thereby facilitating read and write operations in the memory device. The memory cells are typically arranged in a number of rows and columns to provide a memory array. Each memory cell in the array is connected to at least one row or “wordline” and at least one column or “bitline.” Generally speaking, the gate terminal of the access device may be coupled to the wordline while at least one of the remaining terminals (e.g. drain/source) is coupled to the bitline. The other terminal may be coupled to the capacitor. When a voltage is applied to the wordline, the gate of the access device opens and charged particles flow from the bitline to the storage device or vice versa, depending on the mode of operation of the memory cell (e.g. read or write).\nAs can be appreciated, a variety of technologies are used to fabricate the memory cells. Generally speaking, layers of conductive, non-conductive, and semiconductive materials are generally disposed on a substrate to form the access and storage structures described above. Numerous fabrication techniques may be used to facilitate the deposition, masking, and etching steps used to construct the memory cells, as can be ascertained by those skilled in the art. In one technique, a conductive material, such as doped polysilicon, is used to form the wordlines which are used to control the gate of the access devices in the memory cell. Each wordline may include several layers disposed on the conductively doped polysilicon layer, thereby forming a wordline stack.\nDuring formation of the wordline stack, a nitride layer, such as silicon nitride (Si3N4), may be formed on the polysilicon layer to provide an even etch surface. While the nitride layer may be advantageous for part of the fabrication process, it may be the cause of problems in later processing steps. For example, during formation of the wordline stack, a conductive metal material, such as tungsten (W), may be disposed on the polysilicon layer to provide conductive contact to metal layers disposed above the integrated circuit. However, because the thin nitride layer may be used to provide an even etch surface on top of the polysilicon layer, the nitride layer may form a barrier between the conductive metal material and the polysilicon layer. Thus, the nitride layer may ultimately reduce the conductivity between the metal layers and the polysilicon, thereby reducing the effectiveness of the integrated circuits."}
{"text": "Freely dissolved concentration in sediment pore water is a critical measurement that is useful in assessing fate, transport, and bioavailability of hydrophobic pollutants in sediment (Mayer, 2014). Accurate measurement of low aqueous concentrations of hydrophobic compounds is challenging due to the association with colloidal and dissolved organic matter in pore water which has led to the development of passive sampling approaches using well characterized polymeric materials. When the polymeric material is able to reach equilibrium with the sediment pore water, such as in a well-stirred laboratory measurement, the estimation of freely dissolved pore water concentration (Cfree) becomes a trivial exercise based on the known partition constant of the polymeric medium (Mayer, 2003; Ghosh, 2014). However, in several situations, an in-situ measurement in sediment is desired, and such measurements have been challenged by the difficulty in reaching equilibrium between pore water and the polymer as mass transfer through the static depletion layer outside the polymer becomes limiting in the absence of active mixing (Lampert, 2010). It has been shown that for strongly hydrophobic compounds equilibrium may not be achieved in the field even after one year (Lohmann, 2011). Several researchers have adopted the use of performance reference compounds (PRCs) dosed in the polymer to assess the kinetics of mass transfer and correct for non-equilibrium (Huckins, 2006; Huckins, 2002; Fernandez, 2009; Booij, 2010; Booij, 2003; Fernandez, 2014; Oen, 2011; Tomaszewski, 2008; Apell, 2014). While corrections based on PRC loss work reasonably well for compounds with low to midrange hydrophobicity, the corrections become increasingly erroneous for strongly hydrophobic compounds when the departure from equilibrium increases (Apell, 2014). Several approaches for calibration using PRC data have been suggested (Huckins, 2006; Huckins, 2002; Fernandez, 2009; Booij, 2010; Booij, 2003). In all of these approaches, the uncertainties introduced by the PRC correction are larger when the extent of equilibrium is low, which is the case for strongly hydrophobic compounds in the field.\nA primary uncertainty in the PRC correction arises from the fact that nearly always it is the sediment side mass transfer in the immediate vicinity of the passive sampler that controls kinetics and as such, is dependent on the site-specific sorption characteristics of the sediment which can vary across orders of magnitude. For example Hawthorne (Hawthorne, 2006) reported a 3-4 orders of magnitude range for site-specific KocS. Thus, to be able to correct for non-equilibrium and estimate in-situ pore water concentrations we need to first have an estimate of site-specific partitioning of the analytes of interest. The loss of a few PRC compounds and adsorption into the sediment matrix then have to be used to infer the desorption behavior of a large range of analyte compounds from sediment (Fernandez, 2009).\nThinner polymeric materials can be used to increase the surface area to volume ratio and reduce the depletion per unit area. However, even with some of the thinnest polymers practically deployable in the field (e.g., 25 μm thick polyethylene), sediment-side mass transfer limitation can be significant. Making the polymers too thin makes them prone to damage during deployment in sediment, reduces the total mass of polymer sampling material (impacting detection limits), and also poses a physical challenge of insertion in sediments if the area is very large.\nTo address these challenges, the present inventors manipulated the external depletion layer in the sediment side of a passive sampler deployment, which overcomes the slow approach to equilibrium for hydrophobic organic compounds in static sediments. Specifically, the present invention relates to an apparatus and method to mechanically disrupt the static depletion layer outside the polymer surface using periodic vibration performed in-situ. Using the apparatus and method to mechanically disrupt the static depletion layer using periodic vibration, the present invention further relates to a method of determining the freely dissolved concentration of analyte in the porewater of sediment."}
{"text": "Automated clinical chemistry analyzers are well known in the art and are generally used for the automated or semi-automated analysis of patient samples. Typically, prepared patient samples, such as blood, urine, spinal fluid and the like, are placed onto such an analyzer in sample containers such a test tubes. The analyzer pipettes a patient sample and one or more reagents to a reaction cell or cuvette where a chemical analysis of the sample is conducted, usually for a particular analyte of interest, and the results of the analysis are reported.\nAutomated pipettors are employed on such analyzers to transfer the patient samples and reagents as required for the specified analysis. Such pipettors can include a hollow probe having an open end or tip. The probe is, for example, lowered into a sample container that holds a sample, a predetermined volume of sample is withdrawn from the sample container, and the probe is withdrawn from the sample container. The probe is moved, for example, to a position above a reaction cell, is again lowered, and the sample held in the probe is expelled into the reaction cell. Similar actions may be used to pipette and deliver one or more reagents from reagent containers to the reaction cell, either with the same probe or with one or more reagent delivery probes.\nTo prepare such a probe for a subsequent delivery, the probe is washed to eliminate as much as possible any residue from the prior samples or reagents that were handled by the probe. Probe washing may be accomplished by, for example, lowering the probe tip into a wash cell that contains a wash fluid such as water. The wash fluid washes the exterior of the probe tip, and the interior of the probe may be cleaned by aspirating and discharging the wash fluid or, alternatively, discharging a wash fluid through the probe into the wash cell.\nA common problem with probe washing, however, is carryover, that is, the residual fluid or contaminates from a fluid that remain on or in or may be absorbed by the probe despite washing. This residue mixes with subsequent sample or reagents drawn into the probe and can interfere with subsequent analyses.\nAnother problem with probe washing is the time needed to move the probe to a wash station and accomplish the probe washing. Substantial time can be required to wash the probe. For example, if the probe has delivered a sample to a reaction cell, the probe must be raised, moved to a position over a wash cell, and lowered into the cell for washing. Once washing is done, the probe must again be raised and moved on to the next operation.\nStill another drawback of prior systems is washing a probe that has been inadvertently inserted into, for example, a sample or reagent to a depth that conventional washing cups and cells can not clean. Probes often include or are provided externally with some form of liquid level detection. This allows the probe to be inserted until the tip touches the surface of the liquid, where liquid contact is sensed. The probe may then be lowered an additional distance sufficient for the volume of liquid to be drawn into the probe or may be lowered as liquid is drawn into the probe. In either instance, only a relatively short, predetermined portion of the outside of the probe tip is contaminated with the liquid. Thus, the probe tip can be cleaned in a conventional wash cell or cup that is required to clean only a short portion of the probe tip, reducing cleaning time and the volume of washing liquid required.\nIf the liquid level sensing used to detect the contact of the probe tip with the surface of the liquid should fail, however, the probe will be lowered some predetermined distance that may represent the maximum insertion distance permitted to avoid the probe tip striking the bottom of the liquid reservoir. If the vessel actually is full or substantially full of liquid, the entire length of the submerged probe exterior becomes contaminated with the liquid. This length may well exceed the probe length that can be accommodated by the washing cell or cup, requiring operator intervention to manually clean the probe and return the probe, and the automated analyzer in which the probe is used, to service.\nThus, there is a need for a probe washing apparatus and method of use of such an apparatus that overcomes these limitations of the prior art probe washing approaches, including but not limited to reduced carryover, decreased probe washing time, and variable exterior probe tip lengths that may be coated with contaminating liquids."}
{"text": "1. Field of Invention\nThe invention relates to wireless communication systems and a universal address book that can be utilized to manage a variety of user contact information.\n2. Description of the Related Art\nWireless communication systems have developed through various generations, including a first-generation analog wireless phone service (1G), a second-generation (2G) digital wireless phone service (including interim 2.5G and 2.75G networks) and a third-generation (3G) high speed data/Internet-capable wireless service. There are presently many different types of wireless communication systems in use, including Cellular and Personal Communications Service (PCS) systems. Examples of known cellular systems include the cellular Analog Advanced Mobile Phone System (AMPS), and digital cellular systems based on Code Division Multiple Access (CDMA), Frequency Division Multiple Access (FDMA), Time Division Multiple Access (TDMA), the Global System for Mobile access (GSM) variation of TDMA, and newer hybrid digital communication systems using both TDMA and CDMA technologies.\nThe method for providing CDMA mobile communications was standardized in the United States by the Telecommunications Industry Association/Electronic Industries Association in TIA/EIA/IS-95-A entitled “Mobile Station-Base Station Compatibility Standard for Dual-Mode Wideband Spread Spectrum Cellular System,” referred to herein as IS-95. Combined AMPS & CDMA systems are described in TIA/EIA Standard IS-98. Other communications systems are described in the IMT-2000/UM, or International Mobile Telecommunications System 2000/Universal Mobile Telecommunications System, standards covering what are referred to as wideband CDMA (WCDMA), CDMA2000 (such as CDMA2000 1×EV-DO standards, for example) or TD-SCDMA.\nIn wireless communication systems, mobile stations, handsets, or access terminals (AT) receive signals from fixed position base stations (also referred to as cell sites or cells) that support communication links or service within particular geographic regions adjacent to or surrounding the base stations. Base stations provide entry points to an access network (AN)/radio access network (RAN), which is generally a packet data network using standard Internet Engineering Task Force (IETF) based protocols that support methods for differentiating traffic based on Quality of Service (QoS) requirements. Therefore, the base stations generally interact with ATs through an over the air interface and with the AN through Internet Protocol (IP) network data packets.\nIn wireless telecommunication systems, Push-to-talk (PTT) capabilities are becoming popular with service sectors and consumers. PTT can support a “dispatch” voice service that operates over standard commercial wireless infrastructures, such as CDMA, FDMA, TDMA, GSM, etc. In a dispatch model, communication between endpoints (ATs) occurs within virtual groups, wherein the voice of one “talker” is transmitted to one or more “listeners.” A single instance of this type of communication is commonly referred to as a dispatch call, or simply a PTT call. A PTT call is an instantiation of a group, which defines the characteristics of a call. A group in essence is defined by a member list and associated information, such as group name or group identification.\nConventionally, data packets within a wireless communication network have been configured to be sent to a single destination or access terminal. A transmission of data to a single destination is referred to as “unicast”. As mobile communications have increased, the ability to transmit given data concurrently to multiple access terminals has become more important. Accordingly, protocols have been adopted to support concurrent data transmissions of the same packet or message to multiple destinations or target access terminals. A “broadcast” refers to a transmission of data packets to all destinations or access terminals that can be commonly associated such as within a given cell, served by a given service provider, and the like, while a “multicast” refers to a transmission of data packets to a given group of destinations or access terminals. In an example, the given group of destinations or “multicast group” may include more than one and less than all of possible destinations or access terminals that otherwise share a common association. However, it is at least possible in certain situations that the multicast group includes only one access terminal, similar to a unicast, or alternatively that the multicast group comprises all access terminals sharing a common association, similar to a broadcast.\nIn addition to various transmission schemes such as unicast, multicast, broadcast that may be used, a call may generally be classified as a half duplex or a full duplex communication for at least some of the participants. A PTT is a half duplex call that corresponds to a server mediated communication between two or more identified access terminals, regardless of the various configurations used to conduct the PTT calls. Full duplex calls can generally include VoIP calls and other calls that allow simultaneous transmission and reception.\nUsers of wireless devices have many methods of communicating including traditional Voice, Short Message Service (SMS), Multimedia Messaging Service (MMS), Email, and Voice over Internet Protocol/Session Initiation Protocol (VoIP/SIP providers. Furthermore, soon each user may also have PTT/PTX options. Further, each user may have multiple accounts with multiple providers, for example, business and personal use which results in an explosion of per-user contact-methods which are difficult to manage in an integrated fashion on a wireless device. Users may need to be able to manage all of these different per-user contact methods in an integrated fashion from their mobile device or from a web-based application, including the ability to sync with local contact management applications.\nPersonal Information Managers (PIMs) such as, for example, Outlook® by Microsoft Corporation, can synchronize with wireless devices, but often not with a large network provider or global address book. Further, most PIM software does not natively support synching a large number of communication methods to a wireless device in a predictable fashion.\nConventional cellular phones generally have an address book, which can be locally stored or can be stored by the network, such as in the GMS Subscriber Identity Module in a GSM network. The address books of individual cellular phones are typically not synchronized in a useful manner to the address books of other phones. Further, synchronization in the form of a very simple point-to-point provider synchronization exists and can either be synchronized to a network provider such as the address book back up service offered by Verizon Corporation, or to a back-up storage medium. Cross-synchronization amongst users however is not common.\nConventional vCards can be sent from one party to another containing name and address information, phone numbers, URLs, logos, photographs, and even audio clips. vCard is a file format standard for electronic business cards. vCards are often attached to e-mail messages, but can be exchanged in other ways, such as on the World Wide Web. However, vCard data can quickly become outdated/inaccurate and it is not easy to update without sending the card.\nThus, a need exists for managing different per-user contact methods in an integrated fashion from a mobile device or from a web-based application, including the ability to synchronize with local contact management applications. Further, a need exists for providing a universal address book (subject to privacy restrictions) that can provide a source of known-good directory information from the wireless carrier or other provider. For example, permission to access a person's contact record may be determined by a users' relationship in a social networking application such as MySpace or Facebook. Furthermore, a universal address book can allow wireless carriers to inter-operate with each other in a seamless way, providing a much “cleaner” and faster user experience."}
{"text": "Imperative programming languages such as C, C++, Java and Python are the workhorses of software engineering, but by default generate code that runs on only a single CPU core. Pure functional languages such as Haskell can be automatically parallelized, but are difficult for the majority of programmers to use productively, as evidenced by the limited adoption of pure functional languages in production environments. This patent covers a new class of languages known as “lattice-based languages” that is a strict superset of the pure functional paradigm and a strict subset of the imperative paradigm, retaining both the automatic parallelization property of pure functional languages and the practical usefulness of imperative languages.\nIntegrated Development Environments (IDEs) facilitate the editing of program code by providing to a programmer facilities that go beyond direct editing of program source, such as mechanisms for understanding the relationship between parts of a program (e.g. cross-linking a symbol with its definition), transforming a program's source code in useful ways (e.g. allowing for all usages of a variable to be renamed in a single operation), and for integrating the debugging of a program at runtime with the same environment used to edit the source code before compiling and/or running the code. IDEs that display a program as text miss an opportunity to display program structure to a programmer in graphical form, however IDEs that use a purely graphical representation for program editing are difficult for programmers to reason about quickly. This patent covers an IDE that augments the text display of a program in a lattice-based language with a graphical display of data dependencies between expressions in the program.\nFunctional reactive programming has been used in spreadsheets and some programming systems to reduce the amount of computation required to compute a final result by caching or memoizing intermediate values, such that when a part of the spreadsheet or program is changed, only those things that depend upon that changed part (or things that depend on them) need to be recomputed. This patent extends the use of functional reactive programming to lattice-based languages, and covers the use of functional reactive programming in an IDE for lattice-based languages."}
{"text": "Batteries are ubiquitous in modern day technology. Batteries are available in a wide range of types, each specialised for particular applications, from small button cell batteries used in watches, larger Li-ion batteries used in mobile devices, to large industrial scale batteries. There is often a trade-off between the physical size of the battery versus the capacity required for the battery's application, the two of which are directly related.\nAn important subclass of batteries is rechargeable batteries, which provide the functionality of being able to be used repeatedly. Rechargeable batteries are the power source of choice for most portable applications, from cellphones and tablets, to electric motors for hybrid or electric vehicles. However, their limited capacity means that once their charge is expended they need to be connected to an external power source in order to be recharged. For applications such as hybrid or electric vehicles this puts a limit on the range these vehicles can travel before having to be recharged, which may be less than desirable.\nFor cellphones, balancing the desire for batteries to be light and portable against the power drain of large LED displays and power hungry processors, leads to batteries that have to be charged almost daily. The charging process for these batteries is time consuming, typically requiring a few hours for a full charge. It is not always desirable to have to recharge these devices this frequently.\nIn addition, rechargeable battery performance can deteriorate as repeated charge and discharge cycles lead to damage to the battery's electrodes. Causes of such deterioration include sulfation and grid corrosion for lead acid batteries, cell oxidation for Li-ion batteries, and crystalline formations for nickel-based batteries. These result in it becoming harder for the battery to regain its initial charge with the capacity of the battery depleting over time.\nThere is therefore a need to improve the functionality of batteries in order to increase their capacity whilst being able to keep them portable enough for specific applications. In addition, more efficient charge and discharge cycles are desirable in order to keep up with development in electronics."}
{"text": "1. Field of the Invention\nThis invention relates to graphics systems and, more particularly, to graphics systems capable of supporting multiple display streams or channels.\n2. Description of the Related Art\nA computer system typically relies upon its graphics system for producing visual output on the computer screen or display device. Early graphics systems were only responsible for taking what the processor produced as output and displaying it on the screen. In essence, they acted as simple translators or interfaces. Modem graphics systems, however, incorporate graphics processors with a great deal of processing power. They now act more like coprocessors rather than simple translators. This change is due to the recent increase in both the complexity and amount of data being sent to the display device. For example, modern computer displays have many more pixels, greater color depth, and are able to display more complex images with higher refresh rates than earlier models. Similarly, the images displayed are now more complex and may involve advanced techniques such as anti-aliasing and texture mapping.\nAs a result, without considerable processing power in the graphics system, the CPU would spend a great deal of time performing graphics calculations. This could rob the computer system of the processing power needed for performing other tasks associated with program execution and thereby dramatically reduce overall system performance. With a powerful graphics system, however, when the CPU is instructed to draw a box on the screen, the CPU is freed from having to compute the position and color of each pixel. Instead, the CPU may send a request to the video card stating, “draw a box at these coordinates.” The graphics system then draws the box, freeing the processor to perform other tasks.\nGenerally, a graphics system in a computer is a type of video adapter that contains its own processor to boost performance levels. These processors are specialized for computing graphical transformations, so they tend to achieve better results than the general-purpose CPU used by the computer system. In addition, they free up the computer's CPU to execute other commands while the graphics system is handling graphics computations. The popularity of graphics applications, and especially multimedia applications, has made high performance graphics systems a common feature in many new computer systems. Most computer manufacturers now bundle a high performance graphics system with their computing systems.\nIn many applications, it may be useful to have two monitors or displays connected to the same computer system. For example, in some graphical editing applications, it is desirable to use one monitor to show a close-up of an area being edited, while another monitor shows a wider field of view of the object or picture being edited. Alternatively, some users may configure one monitor to display the object being edited and the other monitor to display various palettes or editing options that can be used while editing. Another situation where multiple displays are useful occurs when several users are connected to a single computer. In such a situation, it may be desirable for each user to have their own display. In another situation, it may simply be desirable to have multiple displays that each display a different portion of an image in order to provide a larger display than would otherwise be possible. Another example is stereo goggles, which present different images to their wearer's left and right eyes in order to create a stereo viewing effect. These examples illustrate just a few of the many situations where it is useful to have multiple displays connected to the same computer system.\nGiven the complexity and expense of many graphics systems, it may be desirable to provide a graphics system that can support two displays without duplicating the entire graphics system. Thus, there is a need to be able to share portions of a graphics system between multiple display channels."}
{"text": "“Cloud computing” services provide shared resources, software, and information to computers and other devices upon request. In cloud computing environments, software can be accessible over the Internet rather than installed locally on in-house computer systems. Cloud computing typically involves over-the-Internet provision of dynamically scalable and often virtualized resources. Technological details can be abstracted from the users, who no longer have need for expertise in, or control over, the technology infrastructure “in the cloud” that supports them.\nConventionally, cloud computing services are utilized via computing devices such as desktop computers with internet and web browser capability. Mobile devices such as smart phones and tablets have become popular client devices for accessing and interacting with cloud computing services."}
{"text": "Mobile devices may run applications, commonly known as “apps,” on behalf of their users. These apps may be distributed through one or more app repositories provided by the first-party manufacturer of the device or operating system or through a third party. These apps may engage in network activity on the mobile device, such as through a cellular or Wi-Fi network. Users may install apps, use apps while they are installed, have them installed without using them, and uninstall them."}
{"text": "Invasive cardiac ablation techniques for the treatment of arrhythmias are well known in the art. For example, U.S. Pat. Nos. 5,443,489 and 5,480,422, whose disclosures are incorporated herein by reference, describe systems for ablating cardiac tissue by application of radio-frequency (RF) energy to the tissue through a catheter.\nIn a cardiac ablation procedure, it is important to apply sufficient energy to create a lesion that will block undesired conduction, while minimizing collateral damage to surrounding tissues. Various methods have been proposed for monitoring ablation procedures for this purpose. For example, U.S. Pat. No. 6,743,225, whose disclosure is incorporated herein by reference, proposes to measure electrical activity of the cardiac tissue proximate a lesion site during an ablation treatment, and then to compare the measurements in order to determine whether the lesion is clinically efficacious so as to be able to block myocardial propagation. The electrical activity can include electrical signals corresponding to the local electrogram signal, pacing threshold value, and the like.\nAs another example, U.S. Patent Application Publication 2007/0198007, whose disclosure is incorporated herein by reference, describes methods and devices for monitoring intracardiac ablation progress in near real time, by evaluating capture of a pacing signal while ablation energy is concurrently directed to a target site. Sufficiency of ablation is indicated by failure of signal capture at a maximum predetermined pacing voltage. An electrode in a cardiac catheter is simultaneously used to test pacing capture and to deliver ablation energy."}
{"text": "The present invention relates to communication antennas, and in particular, antennas which have an emitter element positioned in front of a reflector surface.\nCrossed dipole antennas for generating linear or circular polarizations are known. A crossed dipole antenna is known from the article “A wide-band aerial system for circularly polarized waves, suitable for ionospheric research”, G. J. Phillips, IEE Proc., Vol. 98 III, 1951, p. 237-–239. Turnstile antennas are described in various U.S. patents. An example is shown in U.S. Pat. No. 2,086,976 from 1935. The antenna shown includes a mast on which multiple crossed antennas are positioned. There are also numerous textbooks which are concerned with turnstile antennas.\nIn order to improve the directional efficiency, the antenna elements are frequently positioned in front of a metallic reflector surface. This approach is known and is applied in the following two antennas. A dual polarized dipole antenna may be inferred from U.S. Pat. No. 6,313,809 (which essentially corresponds to German Published Application 198 60 121 A1) of the firm Kathrein. These dipole antennas are distinguished in that they include a number of individual dipole elements which are positioned in front of a reflector. The dipole elements are positioned as a dipole square in a top view and each dipole element is fed individually via a symmetrical line.\nA double-polarized multirange dipole antenna may be inferred from U.S. Pat. No. 6,333,720 of the firm Kathrein.\nThe bandwidth of a dipole antenna may be improved by using thick dipoles or bowtie dipole structures. Such a broadband dipole antenna is described in the article entitled “Broad-band half-wave dipole”, M. C. Bailey, IEEE Trans. Antennas Prop., Vol. 32, 1984, p. 410–412. In addition, a broadband antenna having a thick dipole structure is cited in the Antenna Engineering Handbook, R. C. Johnson and H. Jasik, editors, 2nd Edition, McGraw Hill, 1984, on p. 28–11.\nIt is a problem of the antenna arrangements of the prior art, which are used in the field of communication and particularly mobile wireless communication, that the antennas are costly and heavy. This leads to expensive and complicated group antennas.\nWhat is needed therefore, is a broadband dipole antenna which is simple and cost-effective. Further, what is needed is a dipole antenna which is suitable for installation in a group antenna."}
{"text": "In the past, photographs containing multiple images of the same object in motion have been produced by either repetitively firing strobe lamps or repeatedly activating a shutter as the object advanced in front of the camera. However, this technique has serious limitations which have resulted in the photographs being unpredictable, usually containing overlapping and poorly aligned images which defeat the artistic and informational content of the photographs. Furthermore, creating movement with a still object on photographic material has involved the actual movement of a camera assembly through the use of a motorized tripod. Because of the masses involved these systems have proven slow and cumbersome to use. Further, in order to produce good results these systems have required a high degree of percision. in positioning the camera which has been difficult to achieve.\nThe present invention overcomes the foregoing limitations by providing a system which can quickly and accurately move an image with respect to a frame of film and expose different parts of the film to different images of the same object without moving the camera or the object itself."}
{"text": "1. Field of the Invention\nThis invention relates to tension rods, and in particular to tension rods for expanding against opposite walls or other support surfaces for holding curtains, hangers or other articles in place on the tension rod.\n2. Description of the Prior Art\nTension rods incorporating telescoping tubes for being installed between walls having different distances between them are well known in the art. Typical tension rods are known for hanging shower curtains. A number of such rods are known having long coil compression springs which are compressed as the telescoping tubes are pressed together and then released so the spring can expand and press the tubes against the walls. However, those rods show the least amount of tension force when the spring is in its expanded state when advantageously it should have its greatest strength, and show the highest tension force when the spring is in its compressed state. In some cases, threaded inner and outer rods are known. However, such of the rods often require a great number of rotations to get the desired length, and often cannot maintain a tight enough fit over time to maintain the rod in place, essentially under downwardly transverse loads. Such constructions are described in U.S. Pat. No. 704,403 (Thill), U.S. Pat. No. 2,079,267 (Vroom), U.S. Pat. No. 2,919,134 (Zuro) and U.S. Pat. No. 3,521,758 (Guilfoyle, Sr.). In some instances a biasing spring is combined with combinations of threaded rods, such as in U.S. Pat. No. 670,585 (Fowler), U.S. Pat. No. 856,316 (Thurston), U.S. Pat. No. 1,548,053 (Mead), U.S. Pat. No. 2,199,851 (Culver), CH 625601 (Baumann) and JP 2055100 (Wada), but these have not been found effective. Compression spring arrangements in tension rods are very well known, but involve a compromise in being reasonably easy to compress while providing sufficient holding force to keep the tension rod in place between opposing walls while the tension rod is under downward transverse loads. This type of tension rods area disclosed in U.S. Pat. No. 468,987 (Lingley et al.), U.S. Pat. No. 519,840 (Edsall), U.S. Pat. No. 536,272 (Edsall), U.S. Pat. No. 645,543 (Birch), U.S. Pat. No. 647,986 (Roberts), U.S. Pat. No. 841,062 (Snyder et al.), U.S. Pat. No. 961,352 (Walters), U.S. Pat. No. 988,200 (Logsdon), U.S. Pat. No. 1,140,570 (Buckley), U.S. Pat. No. 1,178,994 (Crump), U.S. Pat. No. 1,425,247 (Galbreath), U.S. Pat. No. 1,639,551 (Booth), U.S. Pat. No. 2,032,842 (Gould), U.S. Pat. No. 2,519,996 (Blake), U.S. Pat. No. 2,973,870 (Schoos), U.S. Pat. No. 3,350,120 (Hinrichs), U.S. Pat. No. 3,952,877 (Kindl), U.S. Pat. No. 4,6037,726 (Schweers), U.S. Pat. No. 4,147,199 (Cameron), U.S. Pat. No. 5,242,065 (Hoban) and U.S. Publication Nos. 2008/0156952 (Nathan) and 2009/0101609 (Batshon). A flexible threaded rod without a spring is disclosed in U.S. Pat. No. 5,330,601 (Geltz). Rods are also known which are initially put between opposing walls, and later tightened against those walls. This can be difficult to do where there is no easy access between the walls and may not produce a tight fit. Other tension rods exist where the rods tend to “walk” or slip as outward tension force is applied to the tension rods.\nNone of the foregoing patent disclosures discloses a compression coil spring in a tension rod which maintains nearly constant tension before, during and after the tension rod is installed. No tension rods are known where a compression spring is located in a holding device disposed in one of a pair of telescoping tubes, wherein the tubes are compressed using low manual force, then twisting one of the tubes to tightly lock the compression spring in place, and then releasing the tubes between the walls. The present invention lacks the shortcomings of the prior art, including long compression springs, compression springs whose springs vary from being relatively high when in a compressed state and relatively weak in an expanded state. The present invention in its preferred form is easy for the average person to compress when required for installation, as is not true for the tension rod having a flexible worm gear for providing the tension force rather than a spring. Tension rods according to the preferred embodiments of the present invention do not tend to “walk” or slip along the surface during installation."}
{"text": "The present disclosure relates to an illuminator that emits out illumination light including a laser beam, and a display unit that carries out an image display using such illumination light.\nAn optical module, which is one of main components of a projector (projection-type display unit), is generally composed of an illuminating optical system (illuminator) including a light source and a projecting optical system (projection optical system) including an optical modulation element. In recent years, in the field of such a projector, a small-sized (palm size) and lightweight portable projector called a microprojector has begun to come into wide use. In the past, for this microprojector, an LED (Light Emitting Diode) has been mainly utilized as a light source for an illuminator.\nMeanwhile, a laser has recently drawn attention as a new light source for an illuminator. For example, a projector utilizing a gas laser has been known in the past as a projector using a laser beam of three primary colors including red (R), green (G), and blue (B). As described above, a projector utilizing a laser as a light source is proposed in, for example, Japanese Unexamined Patent Application Publication No. 2005-301164. Use of a laser as a light source allows to achieve a projector with a wide color reproduction range and a reduced power consumption."}
{"text": "1. Field of the Invention\nThe present invention relates to an air-compressing, self-ignition or spark-ignition, four-stroke internal combustion engine having direct fuel injection, turbo-charging, and load-dependent internal exhaust gas recirculation. This exhaust gas recirculation is effected, at least in certain operating ranges, by intervention in the gas change system. Within this gas change system, the mixture formation is essentially controlled by means of the high-speed rotary motion of the fresh air charge, which is produced in an intake channel and is maintained in a combustion chamber having the shape of a solid of revolution.\n2. Description of the Prior Art\nDiesel engines have qualitative governing or control. There is no throttling and, consequently, volume flow is high. Owing to the difficulty of the task of achieving mixture formation immediately before combustion in a minimum of time, provision is made for more air to be drawn into the combustion chambers of the engine than is actually necessary for the combustion of the fuel injected into the combustion chambers. Generally, the greater the proportion of the inducted air, the lower the load of the engine. Moreover, combustion in the low load range is sluggish and takes place at a low temperature level. Therefore, this is the range where the exhaust gases are unfavorable with respect to the unburnt substances; above all, the great amount of oxygen is conducive to the formation of nitric oxides. This is particularly so in the case of direct-injection engines.\nIn recent years, it has become general practice, in turbo-charged internal combustion engines, to return a portion of the exhaust gases to the intake in order to reduce emission of pollutants (nitric oxides and hydro-carbons).\nExhaust gas recirculation reduces the oxygen content of the air for combustion and, consequently, the effective excess air in the fresh gases. In other words, the reaction kinetics of combustion are interfered with by means of the oxygen concentration of the cylinder charge, whereby the combustion process and the exhaust gas composition are influenced.\nAnother very important aspect of exhaust gas recirculation is the reduction of the ignition lag, which includes the time from the start of injection of the fuel to the start of combustion. Ignition lag is a consequence of the higher final compression temperature which results from the higher fresh gas inlet temperature. In addition to other advantages, e.g. the reduction of ignition noise, the shortening of the ignition lag results in an improvement in combustion which, in turn, decreases the emission of pollutants.\nIn an exhaust gas recirculation control system for diesel engines, it is desirable that the amount of air replaced by exhaust gas should be proportional to the amount of surplus air relative to actual air requirements for the combustion of the injected fuel, such that a maximum amount of surplus air is removed from the air flow supplying the engine cylinder without causing unsteady combustion of the fuel in the cylinders. The objective is to achieve maximum efficiency throughout the operating spectrum of the engine with respect to control of nitric oxide emission.\nThe maximum rate of exhaust gas recirculation is required in the lower part load range because this is where the greatest amount of excess air exists. In the full-load range, however, a high rate of exhaust gas circulation would reduce the output which the engine can attain because there is only a small amount of excess air. For this reason, it is necessary for the rate of exhaust gas recirculation to be controlled so that the proportion of recirculated exhaust gas decreases as the load of the engine increases until no exhaust gas at all is recirculated at full load.\nA distinction is made between external and internal exhaust gas recirculation. In the case of external recirculation, the exhaust gases are returned from the exhaust gas ducts via pipes and control devices into the intake channel. In contrast to this, internal exhaust gas recirculation can be implemented in a simpler manner by appropriate intervention in the gas change system. Therefore, internal recirculation offers certain advantages over the external concept; especially, it affords advantages with respect to hydro-carbon emission in the low-load range (due to the hotter exhaust gas). Last but not least, it will also expedite cold starting and warming-up of the engine.\nIntervention in the gas change system in the case of internal exhaust gas recirculation can be effected in various ways (cf. German Auslegeschrift No. 12 22 735 and German Pat. No. 12 42 044), the means employed including, inter alia, double-rise cams for controlling the gas change valves (cf. for instance, German Auslegeschrift No. 14 01 228, German Offenlegungsschrift No. 21 25 368, German Offenlegungsschrift No. 26 38 651, German Offenlegungsschrift No. 27 10 189, German Pat. No. 17 51 473).\nOf the references mentioned, only the German Pat. No. 12 42 044 covers a turbo-charged internal combustion engine.\nThe present invention relates to this type of internal combustion engine, which uses internal exhaust gas recirculation in a similar manner to that already described.\nAn object of the present invention is to improve such an internal combustion engine of the aforementioned general type in such a way that, on the one hand, exhaust gas recirculation does not cause any weakening of the air swirl energy (charge swirl energy) at low loads, which would spoil the desired exhaust gas improvements, and that, on the other hand, assisted cylinder scavenging would be obtained at high loads."}
{"text": "This invention relates to the field of trephine cutting tools of the type used for obtaining samples from humans, animals, and other living or deceased specimens.\nIn the treatment and study of disease and the therapeutic processes and in biomedical studies such as space medicine bone study, it is desirable to collect small specimens of bone from a vertebrate test subject for microscopic and other laboratory study. A guiding and supporting arrangement for conveniently and safely collecting such samples is, for example, shown in the copending patent application, having one common inventor herewith, Ser. No. 06/751,393, filed Jul. 3, 1985 and issued 31 Mar. 1987 as U.S. Pat. No. 4,653,509, and deals principally with accurately positioned sample collection from small animal test subjects.\nThe use of a hollow annular cross-sectioned drill bit having small sawteeth disposed at the working end thereof, i.e., a trephine cutter, for the collection of tissue samples from living test subjects has been known for some time. A frequent problem with the known sample collecting arrangements occurs, however, during the post-sample collection steps wherein the sample is removed from the trephine cutter bit. Heretofore it has been common practice to employ hydraulic pressure or a push rod or some similar forcing apparatus for disengaging the collected sample from the trephine cutter bit. In the course of this disengagement, sufficient pressure is often applied to the collected sample to cause compression damage or sample density alteration as a by-product of the disengagement sequence. For most present-day study purposes, such disturbances of the collected sample are a disabling, sample-destroying consequence. An awareness of such destructive effects occurring during sample collection is, for example, indicated in the article \"Closed Vertebral Biopsy\" by L. S. Fyfe, A. P. J. Henry, and R. C. Julkolland, published in the Journal of Bone and Joint Surgery, Vol 65-B, No. 2, Mar. 1983, wherein the collection of tissue specimens of two millimeters diameter or more is suggested for best accuracy. In the case of bone samples, the tendency of sampled tissue to clog or obstruct the sample collecting apparatus has in particular been a problem with the existing trephine arrangements. The periosteum tissue surrounding bone members in a vertebrate specimen has particular affinity for surgical instruments and thereby materially complicates the sample retrieval process using most existing trephine equipment.\nThe patent art discloses inventive attention to the field of trephine sample collection. This attention is evidenced in the patents of W. Ackerman, U.S. Pat. No. 2,919,692; I. Hallac, U.S. Pat. No. 3,605,721; K. Jamshidi, U.S. Pat. No. 3,628,524; P. W. Hoffsess, U.S. Pat. No. 3,893,445; and L. S. Matthews, U.S. Pat. No. 4,306,570, all of which concern trephine or biopsy sample apparatus providing some arrangement for disengaging the collected sample--arrangements such as a push rod extruder or multiple collecting device diameters. It may be noted, however, that none of these patents is concerned with a trephine sample collecting arrangement wherein the cutter bit can be disassembled to provide convenient access to the sample and also assure the absence of compression disturbances to the sample. The prior art has therefore been unable to provide a fully satisfactory apparatus for collecting high-quality trephine samples conveniently and practically."}
{"text": "Depending on the medical circumstances, reduced pressure may be used for, among other things, reduced-pressure therapy to encourage granulation at a tissue site, draining fluids at a tissue site, closing a wound, reducing edema, promoting perfusion, and fluid management. Common dressings, systems, and methods may be susceptible to leaks and blockage that can cause a reduction in the effectiveness of the therapy or a complete loss of therapy. Such a situation can occur, for example, if the amount of fluid in the dressing or system exceeds the fluid capacity of the dressing or system. Thus, improvements to dressings, systems, and methods that may enhance the management of fluid extracted from a tissue site for increasing the reliability and the useable life of the dressing and system are desirable."}
{"text": "During the drilling of a wellbore, various fluids are typically used in the well for a variety of functions. The fluids may be circulated through a drill pipe and drill bit into the wellbore, and then may subsequently flow upward through wellbore to the surface. During this circulation, the drilling fluid may act to remove drill cuttings from the bottom of the hole to the surface, to suspend cuttings and weighting material when circulation is interrupted, to control subsurface pressures, to maintain the integrity of the wellbore until the well section is cased and cemented, to isolate the fluids from the formation by providing sufficient hydrostatic pressure to prevent the ingress of formation fluids into the wellbore, to cool and lubricate the drill string and bit, and/or to maximize penetration rate.\nIn most rotary drilling procedures the drilling fluid takes the form of a “mud,” i.e., a liquid having solids suspended therein. The solids function to impart desired rheological properties to the drilling fluid and also to increase the density thereof in order to provide a suitable hydrostatic pressure at the bottom of the well. The drilling mud may be either a water-based or an oil-based mud.\nMany wells, especially in oil fields in shale formations (having water sensitivity) and/or deep-water/subsea environments, are drilled with synthetic/oil-based muds or drilling fluids. Because of the extremely high cost of intervention and high production rates, these wells require reliable completion techniques that prevent sand production and maximizes productivity throughout the entire life of the well. One such technique is open-hole gravel packing.\nIn unconsolidated formations, sand control measures are implemented to stabilize formation sand. Common practice for controlling sand displacement includes placement of a gravel pack to hold formation sand in place. The gravel pack is typically deposited around a screen. The gravel pack filters the sand while still allowing formation fluid to flow through the gravel, the screen and a production pipe.\nThere are two principal techniques used for gravel packing open holes: (1) the alternate path technique and (2) alpha-beta packing technique. The latter uses low-viscosity fluids, such as completion brines to carry the gravel from the surface and deposit it into the annulus between a sand-control screen and the wellbore. The alternate path technique, on the other hand, utilizes viscous carrier fluids; therefore the packing mechanisms of these two techniques are significantly different. The alternate path technique allows bypassing of any bridges that may form in the annulus, caused by for example high leakoff into the formation due to filtercake erosion, or exceeding the fracturing pressure, or shale-sloughing/shale-swelling or localized formation collapse on the sand control screens."}
{"text": "Field of the Invention\nThe present invention generally relates to facilitating financial transactions and more particularly to facilitating such transactions with a mobile device.\nRelated Art\nMobile devices, such as cell phones, are being used by more and more people world-wide. In addition to using mobile phones for voice calls, consumers can communicate nearly anytime and anywhere to facilitate information access to mobile services and the Internet. Mobile phones have also become multimedia computing platforms with integral digital cameras for taking pictures and video, playing music, recording conversations, and organizing and planning activities and appointments.\nMore recently, mobile phones have been used in conjunction with on-line payment service providers, such as PayPal, Inc. of San Jose, Calif. With the ever-increasing popularity of the Internet and of Internet commerce, both consumers and sellers are using the Internet to facilitate financial transactions between parties, whether they are individuals or companies. In on-line financial transactions, consumers may use an on-line payment provider to transfer money and receive money over electronic networks, such as the Internet. The money transfer may be for payment of goods or services. The payment providers supply an infra-structure, software, and services that enable users to make and receive payments. Mobile phone users may contract with a payment provider to allow the user to make payments over the phone. Typically, this requires the user to open up an application on the phone, such as through a web browser. The user then accesses his or her on-line account by entering requested information, such as a user name, phone number, email, and/or password. Payment information can then be entered and transmitted to the payment provider, who then transfers funds from the user's account to the payee's account. A confirmation may then be sent to the payer and/or the payee.\nWhile this service gives the consumer flexibility in paying for services anywhere using a mobile phone, the procedure can be cumbersome, time-consuming, and may be prone to fraudulent transfers."}
{"text": "This invention relates generally to a powered stapling device and, more specifically, to a powered stapling device for stapling objects that are in difficult to reach places.\nPowered staple guns serve a variety of purposes and often the structural configuration and operation of the staple gun is customized for a specific purpose. For example, long handled staple guns are used for stapling material on ceilings. Another type of staple gun typically used in construction includes one that operates as a modular powered tool with an interchangeable handle and magazine units that can drive either nails or staples.\nOne type of powered staple gun having a long nose for reaching otherwise inaccessible locations is described in U.S. Pat. No. 3,834,602, to Obergfell (the '602 patent). The '602 patent discloses a powered staple gun with a nosepiece or drive track of substantially increased length that does not require an increased stroke for driving the nail or staple. The powered staple gun is capable of being operated by a pneumatic motor. The staple or nail driven by the powered staple gun of the '602 patent is advanced through a drive track in increments by a series of strokes, which provide the energy for driving the staple or nail. The configuration of the powered staple gun is such that the user must hold the gun substantially perpendicular with respect to a substrate onto which an object is to be stapled. For example, if the user is stapling a linear object, such as cable or wire, the user must hold the gun at a 90 degree angle to the substrate, which results in the staples being driven over the linear object such that the body of the staple is substantially perpendicular to the linear object.\nIt would be desirable to have a powered stapling device that can be used to reach inaccessible or difficult to reach places. In addition, it would be desirable to have a powered stapling device that can drive a staple at a desired angle even though the powered stapling device is aligned with or perpendicular to a linear object that is to be stapled."}
{"text": "Large round bales have become an exceedingly popular package in which to store a variety of crop materials, ranging from alfalfa to corn stalks. The large cylindrical configuration and generally uniform density of the bales enhance their weathering characteristics, allowing the bales to remain in the field until needed. Additionally, the bales concentrate large quantities of crop material in individual packages, reducing the time required to gather and transport a given amount of crop material when needed. Further, other than a tractor operator, no manual labor is required to gather and transport the bales from the field. In sum, the benefits of large round bales are consistent with the continuing trend toward highly mechanized agricultural operations.\nConventional large round balers require discharge of a completed bale before formation of a new bale is started. Such requirement consumes a significant amount of the total baling time.\nSeveral proposed designs exist in the prior art of round balers which appear to be capable of continuous operation. See, for example, U.S. Pat. Nos. 4,035,999 and 4,499,714. However, neither of these proposed designs has achieved commercial acceptance. Apparently, the projected increases in productivity expected to be gained using balers of these designs were outweighed by the complexity, costs, and other unresolved problems associated therewith."}
{"text": "Power converters, such as buck converters, are commonly utilized to convert a high DC voltage to a low DC voltage. A power converter typically includes a high-side switch and a low-side switch connected in a half-bridge configuration. The power converter can include a driver integrated circuit (IC) to control a duty cycle of either or both of the high-side and low-side switches so as to convert a high input voltage to a low output voltage. To improve form factor, performance, and manufacturing cost, it is often desirable to integrate components of a power converter circuit, such as a half-bridge based DC-DC converter or a voltage converter, into a compact power semiconductor package.\nIn a conventional power semiconductor package, individual semiconductor dies are arranged side by side and coupled to a substrate through their corresponding conductive clips, which can undesirably increase electrical resistance and form factor of the power semiconductor package. Also, package design rules to successfully accommodate multiple leadframes and a conductive clip require a large degree of tolerance (i.e. a large clearance space) for manufacturing. Typically, a conductive clip having a leg portion is used to provide sufficient clearance space for necessary electrical connections. However, it is difficult to manufacture the leg portion of the conductive clip to match the exact height of the semiconductor devices in the conventional power semiconductor package. As a result, the leg portion may cause the conductive clip to tilt either toward or away from the semiconductor device, which in turn can cause unreliable electrical connection between the conductive clip and the semiconductor device, and limiting the current carrying capability of the conductive clip. Additionally, the increased package complexity resulting from the use of multiple conductive clips may negatively affect manufacturing time, cost, and package yields.\nThus, there is a need in the art to provide a compact power semiconductor package to with reduced form factor and increased current carrying capability."}
{"text": "The disclosure of the present application relates generally to tissue removal, and more particularly to devices, systems, and methods for percutaneous and mini-invasive valve removal.\nFor decades, percutaneous interventional therapy has been an option for patients with pulmolic, mitral, and/or aortic valvular disease. The treatment preferred in selected patients with pulmolic or mitral stenosis is percutaneous valvuloplasty. According to the current ACC/American Heart Association (AHA) guidelines in patients with calcific aortic stenosis, balloon aortic valvuloplasty (BAV) has been used as a bridge to aortic valve replacement.\nHospital mortality for BAV varies from 3.5% to 13.5%, while serious complications appear in at least 25% of the patients. The durability of BAV is restricted. Consequently, open aortic valve replacement continues to be the best therapy for aortic stenosis in patients who are viable candidates for surgery. The most frequent heart valve operation is the aortic valve replacement. In the United States, from 2% to 7% of individuals older than 65 years suffer from aortic stenosis (AS). The percentage of people over 65 years with aortic stenosis will continue to increase because people are living increasingly longer lives. Aortic stenosis is frequently associated with comorbid risk factors and previous bypass surgery because it is persistent and progressive and occurs frequently in elderly patients. The surgical therapy for AS patients is useful to improve symptoms and prolong life.\nPercutaneous strategies for the treatment of AS began with percutaneous balloon valvuloplasty. Percutaneous balloon valvuloplasty data from the multicenter National Heart, Lung, and Blood Institute (NHLBI) registry shows only mild progress in early hemodynamics, a significant incidence of peripheral vascular complications, a 30-day mortality rate of 7%, and a high incidence of restenosis within 6 months.\nUnsatisfactory results of balloon aortic valvuloplasty (BAV) have led to investigation of percutaneous placement of prosthetic aortic valves. Such devices are being clinically utilized in a small number of cases of high-risk patients. Although percutaneous aortic valve insertion has been performed on extremely high-risk patients, considerable para-valvular leak regurgitation and early mortality have discouraged this approach.\nOne of the biggest problems with percutaneous or transapical aortic valve replacement is the dilatation of a calcific aortic valve prior to delivery of the stent valve device. The consequences of irregular dilation of the aortic valve area are periprosthetic leak, calcium embolization, difficulties with device insertion and possible migration of the device.\nIn order to address these issues, present application discloses a device, system, and method for endovascular resection of the calcific aortic valve having the following features: a valve isolation system, a mechanism to cut and destroy a heavily calcified valve, and a technique and device to prevent migration or embolization of calcific debris into visceral organs. This device will set the stage for better delivery and insertion of the stent valve device."}
{"text": "The need to store digital files, documents, pictures, images and other data continues to increase rapidly. In connection with the electronic storage of data, various systems have been devised for the rapid and secure storage of large amounts of data. Such systems may include a number of storage devices that are used in a coordinated fashion. In particular, data can be distributed across multiple storage devices such that data will not be irretrievably lost if one of the storage devices (or in some cases, more than one storage device) fails. An additional advantage that can be achieved by coordinating operation of a number of individual storage devices is improved data access and/or storage times. Examples of systems that can provide such advantages can be found in the various RAID (redundant array of independent disks) levels that have been developed.\nMany data storage systems, such as RAID array systems, utilize one or more controllers. Such controllers typically provide a read cache and a write cache to improve the performance of the data storage system with respect to input/output (IO) operations. For example, in connection with write operations, a host can provide data for storage that is initially placed in the write cache of a data system controller. Once the data has been placed in the write cache, the controller indicates to the host that the write operation has been completed, even though the data has not yet been stored on a data storage device or devices. This frees the host to move on to other operations while the controller completes the relatively slow process of storing the data that has been placed in the write cache on the storage device or devices.\nIn many applications, the availability of data is a key concern. In particular, in many applications, users rely heavily on the data stored in the data storage system. In these types of applications, unavailability of data stored on the data storage system can result in significant loss of revenue and/or customer satisfaction. Employing a data storage system comprising a RAID array can enhance the availability of the stored data, since if a single disk drive fails, data may still be stored and retrieved from the system.\nThe operation of the data storage system controller (or storage controller) is typically determined through the execution of controller code. In connection with certain storage controllers, multiple copies of controller code can be maintained in storage controller memory. In normal operation, only the code labeled as the primary copy is executed. In order to upgrade code in a typical dual copy storage controller, the primary code that is being replaced is copied from a primary region in memory to a second region in memory. Then, the new, upgraded code is copied to the primary region, thus the new code becomes the primary code. This method of code upgrading works, but it is slow and not robust because it involves a number of time consuming steps and the movement of relatively large amounts of code. In particular, an internal copy of the code from one region to another is required. Because the controller code is typically the size of several megabytes even though the code is compressed, operations copying code from one region and programming it to another region can take many seconds to complete. Furthermore, where there are many pieces of programmable code to copy on the storage controller, it may take many minutes just to complete the copy portion of the upgrade process. The time required becomes even greater in storage controllers that contain more than two sets of controller code. For example, a three copy programmable component, with three regions being labeled as the latest (newest) region, medium region and old region, involves the copy from the medium region to the old region and then from the latest region to the medium region before programming the downloaded code to the latest region. In addition, such methods of code upgrading are error prone because each programmable code component copy operation is an additional step that is subject to failure. The likelihood of such failure increases proportionally with the number of programmable code components the storage controller has and the number of code copies the storage controller maintains. In addition, such methods of code upgrading lack flexibility, as the newly loaded code is always designated as the primary copy or the latest copy designated for execution. Furthermore, the older version of the controller code is always overwritten, leaving the newer copy as the back-up for the downloaded code. This can be undesirable where the newer version of the code has a serious problem. However, storage controllers have been incapable of overwriting any set of code other than the oldest code in connection with downloading upgraded code."}
{"text": "The invention relates to a method and apparatus for searching the synchronization signal of a next frame and thus determining the frame length of the current frame in digital encoded signals. More particularly, the invention relates to a method and apparatus for utilizing a lookup table to determine a search region and backward-searching the synchronization signal of the next frame in the search region.\nVarious digital video/audio processing standards are widely used in the industry to encode the video or audio data into digital encoded signals for further transmission or storing operations. For instance, MPEG audio standard (ISO/IEC 11172-3 and ISO/IEC13818-3) is for compressing and decompressing audio signals with different sampling rates and bit rates. MPEG audio standard defines three layers, which are Layer I, Layer II, and Layer III. The digital encoded signals conforming to MPEG audio standard includes a plurality of frames; each frame includes a plurality of fields. FIG. 1 illustrates the format of a frame specified in the three different layers of the MPEG audio standard. Take MPEG audio layer-1 for example, each frame includes header, alloc, scalefactors, samples, and ANC fields. Though MPEG audio layer I, II, and III have different formats of frames, they have the same header formats, as shown in FIG. 2.\nPEG audio standard defines compressing and decompressing audio signals with different sampling rates and bit rates. For some specific sampling rate and bit rate, frames with varying lengths are used to achieve average bit rate. Since each frame may have different length, the actual length of each frame has to be obtained when the digital encoded signal is to be decoded. In U.S. Pat. No. 5,777,992, a method used to calculate the length of a current frame is disclosed, and the following formula is included:\n      P    =                  Br        N            ×                        n          S                          F          S                      ,where P represents the number of slots included in the current frame; Br is a bit rate of the digital encoded signal; N is the number of bits included in a slot; nS is the number of samples included in a frame; and FS is the sampling rate of the digital encoded signal.\nWhen a decoder of the related art receives a current frame of the above-mentioned digital encoded signal, the decoder first utilizes the above-mentioned formula to calculate the length of the current frame and then perform the following decoding operations toward the current frame. When P calculated by utilizing the formula is an integer, the current frame includes P slot(s); when the calculated P is not an integer, the current frame possibly includes P′r P′+1 slot(s), wherein P′ is the highest integer whose value is less than P. To correctly decode the digital encoded signals with varying frame lengths, the decoder needs an approach to determine the number of slots (equal to either P′ or P′1) included in the current frame.\nIn MPEG audio standard, the header includes a padding bit signaling the length of the frame. If the calculated P by utilizing the above-mentioned formula is not an integer, the decoder of the related art checks the padding bit; if the padding bit is 0, the current frame includes P′ slots; if the padding bit is 1, the current frame includes P′+1 slots. In other words, the decoder of the related art utilizes the above-mentioned formula to calculate P and then determine the length of the current frame by checking the padding bit.\nIn the header of each frame in the digital encoded signal, a synchronization signal is included. The synchronization signal generally owns specific data patterns for easy identification. For digital encoded signals conforming to the MPEG audio standard, the synchronization signal is called ‘syncword’ The data pattern of the syncword is fixed as ‘1111 1111 1111’ in the binary system, which is ‘0xFFF’ in the hexadecimal system. The system of the related art can forward-search the fixed data pattern (‘0xFFF’) from the current frame so as to determine the syncword in a next frame. Afterwards, by calculating the difference of the positions between the syncword of the current frame and the syncword of the next frame, the length of the current frame is obtained.\nThere have been many MPEG audio encoders accumulated in the market, and some does not encode audio signals in a way strictly conforming to MPEG audio standard. According to MPEG standard, the specific data pattern is exclusively reserved for the syncword and should not appear in any other positions in the frames. However, for some digital encoded signals not strictly conforming to MPEG audio standard, the specific data pattern may appear in other positions in the frames. For instance, in MPEG audio standard, ‘0’ should be utilized as a stuffing bit, so the stuffing sequence is ‘0000 . . . ’ However, some encoders utilize ‘1’ as the stuffing bit, so the stuffing sequence is ‘1111 . . . ’ and may be confused with the syncword (‘1111 1111 1111’). Therefore, the above-mentioned forward-searching method does not work due to the possibility that other data (such as the inappropriate stuffing sequence ‘1111 . . . ’ may be incorrectly interpreted as the syncword, and thus an incorrect frame length of the current frame is obtained."}
{"text": "Non-Small Cell Lung Cancer (NSCLC)\nLung cancer has become the number one killer among cancers worldwide (Molina, J. R. et al., Mayo Clin Proc. 2008 May; 83(5): 584-594). Only 15% of all lung cancer patients are alive five (5) years or more after diagnosis (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). The two (2) main types of lung cancer are small cell lung cancer (SCLC) and non-small cell lung cancer (NSCLC), the latter of which accounts for approximately 85% of all cases of lung cancer (Molina, J. R. et al., Mayo Clin Proc. 2008 May; 83(5): 584-594; Navada, S. et al., J Clin Oncol. 2006; 24(18S) suppl: 384S; Sher, T. et al., Mayo Clin Proc. 2008; 83(3): 355-367).\nRisk Factors for NSCLC\nThe primary risk factor for lung cancer is smoking, which accounts for more than 85% of all lung cancer-related deaths (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801; Doll, R. et al., Br Med J 1976; 2: 1525-1536). The risk for lung cancer increases with the number of cigarettes smoked per day and the number of years spent smoking. In addition to the hazard of first-hand smoke, exposed nonsmokers have an increased relative risk for developing lung cancer (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801; Wald N. J. et al., Br Med J 1986; 293: 1217-1222). Radon gas, a radioactive gas that is produced by the decay of radium 226, is the second leading cause of lung cancer. The decay of this isotope leads to the production of substances that emit alpha-particles, which may cause cell damage and therefore increase the potential for malignant transformation (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801; Schrump, D. S. et al., DeVita, Hellman, and Rosenberg's Cancer: Principles & Practice of Oncology, 8th Edition. Vol. 1. Philadelphia: Lippincott Williams & Wilkins; 2008:896-946).\nAsbestos, a mineral compound that breaks into small airborne shards, is a known carcinogen that increases the risk for lung cancer in people exposed to the airborne fibers, especially those who smoke (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). An estimated 3% to 4% of lung cancers are caused by asbestos exposure (Omenn, G. S. et al., Environ Health Perspect 1986; 70: 51-56). Other possible risk factors include recurring lung inflammation, lung scarring secondary to tuberculosis, family history, and exposure to other carcinogens, such as bis(chloromethyl) ether, polycyclic aromatic hydrocarbons, chromium, nickel, and organic arsenic compounds (Fraumeni, J. F., J Natl Cancer Inst 1975; 55: 1039-1046; Janerich D. T. et al., N Engl J Med 1990; 323: 632-636).\nFamilial clustering or aggregation of lung cancer has been reported in the past 60 years, suggesting a hereditary base to disease development (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801; Sellers, T. A. and Yang P., King R A, Rotter J I, Motulsky A G, editors. The Genetic Basis of Common Diseases. 2nd ed. New York, N.Y.: Oxford University Press; 2002. pp. 700-712; Hwang, S. J. et al., Hum Genet. 2003 August; 113(3): 238-243; Bailey-Wilson, J. E. et al., Am J Human Genet. 2004 September; 75(3): 460-474; Hung, R. J. et al., Nature. 2008; 452(7178): 633-637). An increased risk of lung cancer was found in carriers of tumor protein p53 (TP53) germline sequence variations (Hwang, S. J. et al., Hum Genet. 2003 August; 113(3): 238-243). A germline epidermal growth factor receptor (EGFR) T790M sequence variation was reported in a family with multiple cases of NSCLC (Lie, X. and Hemminki, K., Int J Cancer. 2004; 112(3): 451-457). A genome-wide linkage study of 52 extended families identified a major susceptibility locus influencing lung cancer risk at 6q23-25p (Bailey-Wilson, J. E. et al., Am J Human Genet. 2004 September; 75(3): 460-474).\nThree (3) independent genetic studies have found a marker on chromosome 15 associated with lung cancer (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). In all three (3) studies, the risk was approximately 30% higher for people with 1 copy of the marker and 70-80% higher for people with two (2) copies. The region where the marker resides contains 3 genes coding for subunits of the nicotinic acetylcholine receptor, a protein on the cell surface onto which nicotine molecules latch, triggering cell change (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801; Thorgeirsson, T. E. et al., Nature. 2008; 452(7178): 638-642; Hung, R. J. et al., Nature. 2008; 452(7187): 633-637).\nStaging of NSCLC\nThere are five stages (Stage 0 to Stage IV) in NSCLC. Stages I, II and III are further subdivided into A and B subtypes. These stages are assigned based on a Tumor, Node and Metastasis (TMN) staging system (See, cancer staging guidelines, www.NCCN.org). The TMN staging system of NSCLC is summarized in Table 1.\nTABLE 1Stages 0 to IV of NSCLCTMNStage(Tumor, Nodes, Metastasis)DefinitionOccult TX, N0, M0TX = Primary tumor cannot becarcinomaassessed, or tumor proven bythe presence of malignant cellsin sputum or bronchialwashings but not visualized byimaging or bronchoscopy.N0 = No regional lymph nodemetastasis.M0 = No distant metastasis.0Tis, N0, M0Tis = Carcinoma in situ.N0 = No regional lymph nodemetastasis.M0 = No distant metastasis.IAT1a, N0, M0T1a = Tumor = 2 cm inT1b, N0, M0greatest dimension.T1b = Tumor > 2 cm but = 3cm in greatest dimension.N0 = No regional lymph nodemetastasis.M0 = No distant metastasis.IBT2a, N0, M0T2a = Tumor > 3 cm but = 5cm in greatest dimension.N0 = No regional lymph nodemetastasis.M0 = No distant metastasis.IIAT1a, N1, M0T1a = Tumor = 2 cm inT1b, N1, M0greatest dimension.T2a, N1, M0T1b = Tumor > 2 cm but = 3T2b, N0, M0cm in greatest dimension.T2a = Tumor > 3 cm but = 5cm in greatest dimension.T2b = Tumor > 5 cm but = 7cm in greatest dimension.N0 = No regional lymph nodemetastasis.N1 = Metastasis in ipsilateralperibronchial and/or ipsilateralhilar lymph nodes andintrapulmonary nodes,including involvement bydirect extension.M0 = No distant metastasis.IIBT2b, N1, M0T2b = Tumor > 5 cm but = 7T3, N0, M0cm in greatest dimension.T3 = Tumor > 7 cm or one thatdirectly invades any of thefollowing: parietal pleural(PL3) chest wall (includingsuperior sulcus tumors),diaphragm, phrenic nerve,mediastinal pleura or parietalpericardium.Tumor in the main bronchus(<2 cm distal to the carina, butwithout involvement of thecarina.Associated atelectasis orobstructive pneumonitis of theentire lung or separate tumornodule(s) in the same lobe.N0 = No regional lymph nodemetastasis.N1 = Metastasis in ipsilateralperibronchial and/or ipsilateralhilar lymph nodes andintrapulmonary nodes,including involvement bydirect extension.M0 = No distant metastasis.IIIAT1a, N2, M0T1a = Tumor = 2 cm inT1b, N2, M0 greatest dimension.T2a, N2, M0 T1b = Tumor > 2 cm but = 3T2b, N2, M0cm in greatest dimension.T3, N1, M0T2a = Tumor > 3 cm but = 5T3, N2, M0cm in greatest dimension.T4, N0, M0T2b = Tumor > 5 cm but = 7T4, N1, M0cm in greatest dimension.T3 = Tumor > 7 cm or one thatdirectly invades any of thefollowing: parietal pleural(PL3) chest wall (includingsuperior sulcus tumors),diaphragm, phrenic nerve,mediastinal pleura or parietalpericardium.Tumor in the main bronchus(<2 cm distal to the carina butwithout involvement of thecarina).Associated atelectasis orobstructive pneumonitis of theentire lung or separate tumornodule(s) in the same lobe.T4 = Tumor of any size thatinvades any of the following:mediastinum, heart, greatvessels, trachea, recurrentlaryngeal nerve, esophagus,vertebral body, carina orseparate tumor nodule(s) in adifferent ipsilateral lobe.N0 = No regional lymph nodemetastasis.N1 = Metastasis in ipsilateralperibronchial and/or ipsilateralhilar lymph nodes andintrapulmonary nodes,including involvement bydirect extension.N2 = Metastasis in ipsilateralmediastinal and/or subcarinallymph node(s).M0 = No distant metastasis.IIIBT1a, N3, M0T1a = Tumor = 2 cm inT1b, N3, M0greatest dimension.T2a, N3, M0T1b = Tumor > 2 cm but = 3T2b, N3, M0cm in greatest dimension.T3, N3, M0T2a = Tumor > 3 cm but = 5T4, N2, M0cm in greatest dimension.T4, N3, M0T2b = Tumor > 5 cm but = 7cm in greatest dimension.T3 = Tumor > 7 cm or one thatdirectly invades any of thefollowing: parietal pleural(PL3) chest wall (includingsuperior sulcus tumors),diaphragm, phrenic nerve,mediastinal pleura or parietalpericardium.Tumor in the main bronchus(<2 cm distal to the carina butwithout involvement of thecarina).Associated atelectasis orobstructive pneumonitis of theentire lung or separate tumornodule(s) in the same lobe.T4 = Tumor of any size thatinvades any of the following:mediastinum, heart, greatvessels, trachea, recurrentlaryngeal nerve, esophagus,vertebral body, carina orseparate tumor nodule(s) in adifferent ipsilateral lobe.N2 = Metastasis in ipsilateralmediastinal and/or subcarinallymph node(s).N3 = Metastasis incontralateral mediastinal,contralateral hilar, ipsilateralor contralateral scalene orsupraclavicular lymphnode(s).M0 = No distant metastasis.IVAny T, any N, M1ATX = Primary tumor cannot beAny T, any N, M1bassessed, or tumor proven bythe presence of malignant cellsin sputum or bronchialwashings but not visualized byimaging or bronchoscopy.T0 = No evidence of primarytumor.Tis = Carcinoma in situ.T1 = Tumor = 3 cm in greatestdimension, surrounded bylung or visceral pleura,without bronchoscopicevidence of invasion moreproximal than the lobarbronchus (i.e., not in the mainbronchus).T1a = Tumor = 2 cm ingreatest dimension.T1b = Tumor > 2 cm but = 3cm in greatest dimension.T2 = Tumor > 3 cm but = 7 cmor tumor with any of thefollowing features (T2 tumorswith these features areclassified T2a if = 5 cm):involves main bronchus, = 2distal to the carina, invadesvisceral pleura (PL1 or PL2)or is associated withatelectasis or obstructivepneumonitis that extends tothe hilar region but does notinvolve the entire lung.T2a = Tumor > 3 cm but = 5cm in greatest dimension.T2b = Tumor > 5 cm but = 7cm in greatest dimension.T3 = Tumor > 7 cm or one thatdirectly invades any of thefollowing: parietal pleura(PL3) chest wall (includingsuperior sulcus tumors),diaphragm, phrenic nerve,mediastinal pleura or parietalpericardium.Tumor in the main bronchus(<2 cm distal to the carina butwithout involvement of thecarina).Associated atelectasis orobstructive pneumonitis of theentire lung or separate tumornodule(s) in the same lobe.T4 = Tumor of any size thatinvades any of the following:mediastinum, heart, greatvessels, trachea, recurrentlaryngeal nerve, esophagus,vertebral body, carina orseparate tumor nodule(s) in adifferent ipsilateral lobe.NX = Regional lymph nodescannot be assessed.N0 = No regional lymph nodemetastasis.N1 = Metastasis in ipsilateralperibronchial and/or ipsilateralhilar lymph nodes andintrapulmonary nodes,including involvement bydirect extension.N2 = Metastasis in ipsilateralmediastinal and/or subcarinallymph node(s).N3 = Metastasis incontralateral mediastinal,contralateral hilar, ipsilateralor contralateral scalene orsupraclavicular lymphnode(s).M0 = No distant metastasis.M1 = Distant metastasis.M1a = Separate tumornodule(s) in a contralaterallobe tumor with pleuralnodules or malignant pleural(or pericardial) effusion.M1b = Distant metastasis.\nScreening for NSCLC\nComputed Tomography (CT)\nComputed tomography (CT) enables small lesions (in the region of 3 cm) to be detected (Read, C. et al., Primary Care Respiratory Journal (2006) 15, 332-336). Until recently, the slow scanning times and high radiation exposure of CT made it an unsuitable tool for screening. However, with the development and refinement of spiral CT, it has become possible to significantly reduce scanning times and radiation exposure. Results of CT screening have been encouraging, with approximately 80-90% of CT screen-detected bronchial carcinomasdiagnosed as early stage 1 tumors (Kaneko, M. et al., Radiology 1996; 201: 798-802; Swensen, S. J. et al., Am J Respir Crit Care Med 2002; 165: 508-513; Henschke, C. I. et al., Radiology 2006; 239(2): 586-590).\nSputum Cytometry\nThe diagnostic yield of sputum cytology is known to vary in relation to tumor location (Read, C. et al., Primary Care Respiratory Journal (2006) 15, 332-336). It has greatest use in the identification of central tumors and is of little or no value in the identification of peripheral cancers (Thunnissen F. B. J. M., J Clin Pathol 2003; 56: 805-810). However, it has been suggested that its sensitivity can be improved through the use of molecular genetic and immunocytochemical markers of malignancy (Marek, W. et al., Eur Respir J 2001; 18: 942-950; Xing, S. et al., Eur Respir J 2005; 25: 956-963).\nAutofluorescence Bronchoscopy\nAutofluorescence bronchoscopy has a high sensitivity for detection of abnormal cells (Read, C. et al., Primary Care Respiratory Journal (2006) 15, 332-336). This technique exploits the differences in the fluorescence properties of bronchial mucosa compared to mucosa of pre-invasive and invasive disease.\nThe best known instrument for autofluorescence bronchoscopy is the Lung Imaging Fluorescence Endoscopy (LIFE) system (Lam, S. et al., J Thorac Cardiovasc Surg 1993; 105: 1035-1040). This instrument uses a blue (442 nm) helium cadmium laser to illuminate bronchial mucosa, and the resulting fluorescence is then digitized into a real-time video image. Other devices, such as the D-light auto-fluorescence system, do not require image-intensifying cameras, but use optical filters incorporated into the bronchoscope (Leonhard, M., Diagn Therap Endosc 1999; 5: 71-75).\nThe majority of published studies on autofluorescence bronchoscopy have shown a significant increase in diagnostic sensitivity for dysplasia and carcinoma in situ detection (Read, C. et al., Primary Care Respiratory Journal (2006) 15, 332-3; Lam, S. et al., Chest 1998; 113: 696-702; Vermylen, P. et al., Lung Cancer 1999; 25: 161-168; Venmans B. J. et al., Diagn Therap Endosc 1999; 5: 77-84; Ikeda, N. et al., Diagn Therap Endosc 1997; 3: 197-201; Yokomise, H. et al., J Bronchol 1997; 4: 205-208; Hirsch, F. R. et al., J Natl Cancer Inst 2001; 93: 1385-1391).\nBiomarkers\nSeveral biomarkers have emerged as prognostic and predictive markers for NSCLC (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). A prognostic biomarker, which is an indicator of innate tumor aggressiveness, is a biomolecule that indicates patient survival independent of the treatment received. A predictive biomarker is a biomolecule that indicates therapeutic efficacy, i.e., an interaction exists between the biomolecule and therapy that impacts patient outcome. Among these biomarkers, evidence is strongest for EGFR, the 5′ endonuclease of the nucleotide excision repair complex (ERCC1), Kirsten rat sarcoma viral oncogene homolog (K-ras), and the regulatory subunit of ribonucleotide reductase (RRM1).\nThe presence of the epidermal growth factor receptor (EGFR) exon 19 deletion or exon 21 L858R mutation does not seem to be prognostic of survival for patients with NSCLC, independent of therapy (Tsao, M. S. et al., N Engl J Med 2005; 353: 133-144). However, the presence of the EGFR exon 19 deletion or exon 21 L858R mutation is predictive of treatment benefit from EGFR-TKI therapy (Sequist, L. V. et al., J Clin Oncol 2008; 26: 2442-2449; Miller, V. A., J Clin Oncol 2008; 26: 1472-1478). High ERCC1 levels are prognostic of better survival for patients with NSCLC when compared with low levels of ERCC1 expression, independent of therapy (Simon, G. R. et al., Chest 2005; 127: 978-983; Olaussen K. A. et al., N Engl J Med 2006; 355: 983-991). High levels of ERCC1 expression are also predictive of poor response to platinum-based chemotherapy (Olaussen K. A. et al., N Engl J Med 2006; 355: 983-991; Bepler, G. et al., J Clin Oncol 2006; 24: 4731-4737). The presence of K-ras mutations is prognostic of poor survival for patients with NSCLC when compared with the absence of these mutations, independent of therapy (Slebos, R. J. et al., N Engl J Med 1990; 323: 561-565). Presence of K-ras mutations is also predictive of lack of benefit from platinum/vinorelbine chemotherapy or EGFR-TKI therapy (Miller, V. A., J Clin Oncol 2008; 26: 1472-1478; Tsao, M. S. et al., J Clin Oncol 2007; 25: 5240-5247). High RRM1 levels are prognostic of better survival for patients with NSCLC compared with low levels of RRM1 expression, independent of therapy (Bepler, G. et al., J Clin Oncol 2004; 22: 1878-1885; Zheng, Z. et al., N Engl J Med 2007; 356: 800-808). High levels of RRM1 expression are also predictive of poor response to gemcitabine-based chemotherapy (Bepler, G. et al., J Clin Oncol 2008; 26(Suppl 1): Abstract 8033; Reynolds, C. et al., J Clin Oncol 2009; 27: 5808-5815).\nTreatment of NSCLC\nSurgery\nIn general, for patients with stage I or II disease, surgery provides the best chance for cure (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). The surgical procedure used depends on the extent of disease and the cardiopulmonary reserve of the patient. If anatomically appropriate and margin-negative resection can be achieved, lung-sparing anatomic resection (sleeve lobectomy) is preferred over pneumonectomy. Otherwise, lobectomy or pneumonectomy are performed if physiologically feasible (Boffa, D. J. et al., J Thorac Cardiovasc Surg 2008; 135: 247-254; Scott, W. J. et al., Chest 2007; 132: 234S-242S). Resection (including wedge resection) is preferred over ablation (i.e., radiofrequency ablation [RFA], cryotherapy, stereotactic radiation) (Scott, W. J. et al., Chest 2007; 132: 234S-242S). However, whether lung-sparing surgeries (i.e., sublobular resection), such as segmentectomy and wedge resection, are useful in patients with severely reduced pulmonary function who are otherwise not candidates for surgery is controversial (Scott, W. J. et al., Chest 2007; 132: 234S-242S; Ginsberg, R. J. et al., Ann Thorac Surg 1995; 60: 615-622; Koike, T. et al., J Thorac Cardiovasc Surg 2003; 125: 924-928).\nRadiation Therapy (RT)\nRT can be used as 1) an adjunct for patients with resectable NSCLC who have no contraindications for surgery; 2) the primary local treatment (i.e., definitive RT) for patients with medically inoperable or unresectable NSCLC; and/or 3) an important palliative modality for patients with incurable NSCLC (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). However, conventional RT alone only results in a median survival of 10 months and a 5-year survival rate of 5% (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801).\nChemotherapy\nMany drugs are active against stage IV NSCLC. These drugs include the taxanes (paclitaxel, docetaxel), vinorelbine, etoposide, pemetrexed, the camptothecin analogs (irinotecan), and gemcitabine. Combinations using many of these drugs produce 1-year survival rates of 30% to 40% and are superior to single agents. Regimens include carboplatin/paclitaxel, cisplatin/paclitaxel, cisplatin/vinorelbine, gemcitabine/cisplatin, cisplatin/pemetrexed, and docetaxel/cisplatin (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801; Ohe, Y. et al., Ann Oncol 2007; 18: 317-323; Fossella, F. et al., J Clin Oncol 2003; 21: 3016-3024; Smit, E. F. et al., J Clin Oncol 2003; 21: 3909-3917; Zatloukal, P. et al., Lung Cancer 2004; 46: 87-98; Scagliotti, G. V. et al., J Clin Oncol 2008; 26: 3543-3551).\nAlthough chemotherapy is appropriate for many patients with lung cancer, there is a consensus that the use of traditional chemotherapeutic agents to treat NSCLC has reached a therapeutic plateau (Molina, J. R. et al., Mayo Clin Proc 2008 May; 83(5): 584-594).\nTargeted Therapy\nSpecific targeted therapies have been developed for treating advanced lung cancer (Sandler, A. B. et al., Clin Cancer Res 2004; 10: 4258s-4262s; Giaccone, G. et al., J Clin Oncol 2005; 23: 3235-3242). Bevacizumab is a recombinant monoclonal antibody that blocksvascular endothelial growth factor (VEGF). Erlotinib is a small molecule inhibitor of EGFR. Cetuximab is a monoclonal antibody that targets EGFR.\nIn 2006, the FDA approved bevacizumab for patients with unresectable, locally advanced, recurrent, or metastatic nonsquamous NSCLC (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). It has been recommended that bevacizumab be used in combination with paclitaxel and carboplatin for select patients with advanced nonsquamous NSCLC (Sandler, A. et al., N Engl J Med 2006; 355: 2542-2550). To undergo treatment with bevacizumab and chemotherapy, patients must have nonsquamous NSCLC and no history of hemoptysis. Any regimen with a high risk for thrombocytopenia (i.e., possible bleeding) should be used with caution when combined with bevacizumab (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801).\nIn 2004, the FDA approved Erlotinib for the treatment of patients with locally advanced or metastatic NSCLC after failure of at least one prior chemotherapy regimen (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). However, erlotinib can also be given as first-line therapy in patients with advanced or metastatic NSCLC who have a known active EGFR mutation or gene amplification (Sequist, L. V. et al., Oncologist 2007; 12: 90-98; N Eng J Med 2009; 361: 947-957; Inoue, A. et al., J Clin Oncol 2009; 27: 1394-1400).\nA recent phase III randomized trial assessed cisplatin/vinorelbine with or without cetuximab for patients with advanced NSCLC (Pirker, R. et al., Lancet 2009; 373: 1525-1531). The results indicated that adding cetuximab slightly increased overall survival (11.3 vs. 10.1 months; P=0.04) (Pirker, R. et al., Lancet 2009; 373: 1525-1531; Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801).\nHowever, there is accumulating evidence that NSCLC acquires resistance to these specific targeted therapies (e.g., erlotinib and gefitinib) and how cancers such as NSCLC become resistant to, for example, EGFR inhibitors (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801). The most common known mechanism is the acquisition of a secondary mutation in EGFR, known as T790M, which renders the kinase resistant to erlotinib and gefitinib (Nguyen, K. S. et al., Clin Lung Cancer 2009; 10: 281-289; Gazdar, A. F., Oncogene 2009; 28(Suppl 1): S24-S31). Amplification of the MET oncogene is another validated resistance mechanism (Ettinger, D. S. et al., J Natl Compr Canc Netw 2010; 8: 740-801).\nDNA Damage-Signaling in Human Tumors\nDNA damage signaling and checkpoint control pathways are among the most commonly mutated networks in human tumors (Morandell, S. et al., Cell Reports 5, 868-877, Nov. 27, 2013; Negrini, S. et al., Nat Rev Mol Cell Biol 2010; 11: 220-228). The three major DNA damage-responsive cell cycle checkpoints are the G1/S checkpoint, intra S-phase checkpoint, and the G2/M checkpoint. In response to DNA damage, eukaryotic cells activate complex signaling networks that mediate DNA repair and cell cycle arrest or, if the damage is extensive, trigger apoptosis (Ciccia, A. and Elledge, S. J., Mol Cell 2010; 40: 179-204). Three canonical protein kinase pathways in both normal and cancer cells arrest the cell cycle in response to damaged DNA: the Ataxia-Telangiectasia and Rad-3 related through Chk1 (ATR-Chk1) pathway, the Ataxia-Telangiectasia mutated through Chk2 (ATM-Chk2) pathway, and the stress-activated protein kinases p38 mitogen-activated protein kinase (MAPK) and its substrate MAPKAP kinase-2 (MK2) (Morandell, S. et al., Cell Reports 5, 868-877, Nov. 27, 2013). The ATM/Chk2 pathway responds primarily to DNA double strand breaks, while the ATRpChk1 pathway is activated by bulky DNA lesions, and following replication fork collapse during S-phase. In contrast to the DNA damage-specific activation of Chk1 and Chk2, the p38MAPK pathway is a general stress-activated kinase pathway that responds to various cellular stimuli, including cytokines, hyperosmolarity, and UV irradiation.\nTumor suppressor protein p53 is a major downstream effector of the aforementioned DNA damage kinase pathways. In normal cells, p53-dependent signaling results in G1 arrest, mainly mediated by transcriptional upregulation of p21. P21 also appears to play a role in sustaining the G2 checkpoint after gamma-irradiation. If DNA damage is extensive, however, p53-dependent pathways target the damaged cell for apoptotic cell death. The MK2 pathway is critical for arresting the cell cycle after genotoxic stress, including cisplatin-induced DNA crosslinks and topoisomerase-inhibitor-induced DNA strand breaks only in tumor cells that lack functional p53 (Manke, I. A. et al., Mol Cell 2005; 17: 37-48; Reinhardt, H. C. et al., Cancer Cell 2007; 11: 175-189). Both the ATRChk1 pathway and the p38-MK2 pathway are required for effective cell-cycle checkpoint function in the absence of p53 (Reinhardt, H. C. et al., Mol Cell 2010; 40: 34-49).\nRecently, Morandell et al. (Cell Reports 5, 868-877, Nov. 27, 2013) showed that, in response to genotoxic chemotherapy, MK2 is essential for the survival of NSCLC tumor cells that lack functional p53 but is dispensable in p53-proficient cells. MK2 was found to specifically sensitize p53-deficient tumors to the DNA-damaging agent cisplatin but had no effect on the treatment response of p53-proficient cancer cells. This suggests a potential for enhancement chemosensitization of p53-deficient tumors to DNA-damaging chemotherapy in vivo through synthetic lethality between p53 and MK2 (Morandell, S. et al., Cell Reports 5, 868-877, Nov. 27, 2013).\nKinases\nKinases are a ubiquitous group of enzymes that catalyze the phosphoryl transfer reaction from a phosphate donor (usually adenosine-5′-triphosphate (ATP)) to a receptor substrate. A classification of all available kinase sequences (approximately 60,000 sequences) indicates kinases can be grouped into 25 families of homologous (meaning derived from a common ancestor) proteins. These kinase families are assembled into 12 fold groups based on similarity of structural fold. Further, 22 of the 25 families (approximately 98.8% of all sequences) belong to 10 fold groups for which the structural fold is known. Of the other 3 families, polyphosphate kinase forms a distinct fold group, and the 2 remaining families are both integral membrane kinases and comprise the final fold group. These fold groups not only include some of the most widely spread protein folds, such as Rossmann-like fold (three or more parallel β strands linked by two α helices in the topological order β-α-β-α-β), ferredoxin-like fold (a common α-β protein fold with a signature βαββαβ secondary structure along its backbone), TIM-barrel fold (meaning a conserved protein fold consisting of eight α-helices and eight parallel β-strands that alternate along the peptide backbone), and antiparallel β-barrel fold (a beta barrel is a large beta-sheet that twists and coils to form a closed structure in which the first strand is hydrogen bonded to the last), but also all major classes (all α, all β, α+β, α/β) of protein structures. Within a fold group, the core of the nucleotide-binding domain of each family has the same architecture, and the topology of the protein core is either identical or related by circular permutation. Homology between the families within a fold group is not implied.\nGroup I (23,124 sequences) kinases incorporate protein S/T-Y kinase, atypical protein kinase, lipid kinase, and ATP grasp enzymes and further comprise the protein S/T-Y kinase, and atypical protein kinase family (22,074 sequences). These kinases include: choline kinase (EC 2.7.1.32); protein kinase (EC 2.7.137); phosphorylase kinase (EC 2.7.1.38); homoserine kinase (EC 2.7.1.39); I-phosphatidylinositol 4-kinase (EC 2.7.1.67); streptomycin 6-kinase (EC 2.7.1.72); ethanolamine kinase (EC 2.7.1.82); streptomycin 3′-kinase (EC 2.7.1.87); kanamycin kinase (EC 2.7.1.95); 5-methylthioribose kinase (EC 2.7.1.100); viomycin kinase (EC 2.7.1.103); [hydroxymethylglutaryl-CoA reductase (NADPH2)] kinase (EC 2.7.1.109); protein-tyrosine kinase (EC 2.7.1.112); [isocitrate dehydrogenase (NADP+)] kinase (EC 2.7.1.116); [myosin light-chain] kinase (EC 2.7.1.117); hygromycin-B kinase (EC 2.7.1.119); calcium/calmodulin-dependent protein kinase (EC 2.7.1.123); rhodopsin kinase (EC 2.7.1.125); [beta-adrenergic-receptor] kinase (EC 2.7.1.126); [myosin heavy-chain] kinase (EC 2.7.1.129); [Tau protein] kinase (EC 2.7.1.135); macrolide 2′-kinase (EC 2.7.1.136); I-phosphatidylinositol 3-kinase (EC 2.7.1.137); [RNA-polymerase]-subunit kinase (EC 2.7.1.141); phosphatidylinositol-4,5-bisphosphate 3-kinase (EC 2.7.1.153); and phosphatidylinositol-4-phosphate 3-kinase (EC 2.7.1.154). Group I further comprises the lipid kinase family (321 sequences). These kinases include: I-phosphatidylinositol-4-phosphate 5-kinase (EC 2.7.1.68); I D-myo-inositol-triphosphate 3-kinase (EC 2.7.1.127); inositol-tetrakisphosphate 5-kinase (EC 2.7.1.140); I-phosphatidylinositol-5-phosphate 4-kinase (EC 2.7.1.149); I-phosphatidylinositol-3-phosphate 5-kinase (EC 2.7.1.150); inositol-polyphosphate multikinase (EC 2.7.1.151); and inositol-hexakiphosphate kinase (EC 2.7.4.21). Group I further comprises the ATP-grasp kinases (729 sequences) which include inositol-tetrakisphosphate I-kinase (EC 2.7.1.134); pyruvate, phosphate dikinase (EC 2.7.9.1); and pyruvate, water dikinase (EC 2.7.9.2).\nGroup II (17,071 sequences) kinases incorporate the Rossman-like kinases. Group II comprises the P-loop kinase family (7,732 sequences). These include gluconokinase (EC 2.7.1.12); phosphoribulokinase (EC 2.7.1.19); thymidine kinase (EC 2.7.1.21); ribosylnicotinamide kinase (EC 2.7.1.22); dephospho-CoA kinase (EC 2.7.1.24); adenylylsulfate kinase (EC 2.7.1.25); pantothenate kinase (EC 2.7.1.33); protein kinase (bacterial) (EC 2.7.1.37); uridine kinase (EC 2.7.1.48); shikimate kinase (EC 2.7.1.71); deoxycytidine kinase (EC 2.7.1.74); deoxyadenosine kinase (EC 2.7.1.76); polynucleotide 5′-hydroxyl-kinase (EC 2.7.1.78); 6-phosphofructo-2-kinase (EC 2.7.1.105); deoxyguanosine kinase (EC 2.7.1.113); tetraacyldisaccharide 4′-kinase (EC 2.7.1.130); deoxynucleoside kinase (EC 2.7.1.145); adenosylcobinamide kinase (EC 2.7.1.156); polyphosphate kinase (EC 2.7.4.1); phosphomevalonate kinase (EC 2.7.4.2); adenylate kinase (EC 2.7.4.3); nucleoside-phosphate kinase (EC 2.7.4.4); guanylate kinase (EC 2.7.4.8); thymidylate kinase (EC 2.7.4.9); nucleoside-triphosphate-adenylate kinase (EC 2.7.4.10); (deoxy)nucleoside-phosphate kinase (EC 2.7.4.13); cytidylate kinase (EC 2.7.4.14); and uridylate kinase (EC 2.7.4.22). Group II further comprises the phosphoenolpyruvate carboxykinase family (815 sequences). These enzymes include protein kinase (HPr kinase/phosphatase) (EC 2.7.1.37); phosphoenolpyruvate carboxykinase (GTP) (EC 4.1.1.32); and phosphoenolpyruvate carboxykinase (ATP) (EC 4.1.1.49). Group II further comprises the phosphoglycerate kinase (1,351 sequences) family. These enzymes include phosphoglycerate kinase (EC 2.7.2.3) and phosphoglycerate kinase (GTP) (EC 2.7.2.10). Group II further comprises the aspartokinase family (2,171 sequences). These enzymes include carbamate kinase (EC 2.7.2.2); aspartate kinase (EC 2.7.2.4); acetylglutamate kinase (EC 2.7.2.8 1); glutamate 5-kinase (EC 2.7.2.1) and uridylate kinase (EC 2.7.4.). Group II further comprises the phosphofructokinase-like kinase family (1,998 sequences). These enzymes include 6-phosphofrutokinase (EC 2.7.1.11); NAD(+) kinase (EC 2.7.1.23); I-phosphofructokinase (EC 2.7.1.56); diphosphate-fructose-6-phosphate I-phosphotransferase (EC 2.7.1.90); sphinganine kinase (EC 2.7.1.91); diacylglycerol kinase (EC 2.7.1.107); and ceramide kinase (EC 2.7.1.138). Group II further comprises the ribokinase-like family (2,722 sequences). These enzymes include: glucokinase (EC 2.7.1.2); ketohexokinase (EC 2.7.1.3); fructokinase (EC 2.7.1.4); 6-phosphofructokinase (EC 2.7.1.11); ribokinase (EC 2.7.1.15); adenosine kinase (EC 2.7.1.20); pyridoxal kinase (EC 2.7.1.35); 2-dehydro-3-deoxygluconokinase (EC 2.7.1.45); hydroxymethylpyrimidine kinase (EC 2.7.1.49); hydroxyethylthiazole kinase (EC 2.7.1.50); I-phosphofructokinase (EC 2.7.1.56); inosine kinase (EC 2.7.1.73); 5-dehydro-2-deoxygluconokinase (EC 2.7.1.92); tagatose-6-phosphate kinase (EC 2.7.1.144); ADP-dependent phosphofructokinase (EC 2.7.1.146); ADP-dependent glucokinase (EC 2.7.1.147); and phosphomethylpyrimidine kinase (EC 2.7.4.7). Group II further comprises the thiamin pyrophosphokinase family (175 sequences) which includes thiamin pyrophosphokinase (EC 2.7.6.2). Group II further comprises the glycerate kinase family (107 sequences) which includes glycerate kinase (EC 2.7.1.31).\nGroup III kinases (10,973 sequences) comprise the ferredoxin-like fold kinases. Group III further comprises the nucleoside-diphosphate kinase family (923 sequences). These enzymes include nucleoside-diphosphate kinase (EC 2.7.4.6). Group III further comprises the HPPK kinase family (609 sequences). These enzymes include 2-amino-4-hydroxy-6-hydroxymethyldihydropteridine pyrophosphokinase (EC 2.7.6.3). Group III further comprises the guanido kinase family (324 sequences). These enzymes include guanidoacetate kinase (EC 2.7.3.1); creatine kinase (EC 2.7.3.2); arginine kinase (EC 2.7.3.3); and lombricine kinase (EC 2.7.3.5). Group III further comprises the histidine kinase family (9,117 sequences). These enzymes include protein kinase (histidine kinase) (EC 2.7.1.37); [pyruvate dehydrogenase (lipoamide)] kinase (EC 2.7.1.99); and [3-methyl-2-oxybutanoate dehydrogenase(lipoamide)] kinase (EC 2.7.1.115).\nGroup IV kinases (2,768 sequences) incorporate ribonuclease H-like kinases. These enzymes include hexokinase (EC 2.7.1.1); glucokinase (EC 2.7.1.2); fructokinase (EC 2.7.1.4); rhamnulokinase (EC 2.7.1.5); mannokinase (EC 2.7.1.7); gluconokinase (EC 2.7.1.12); L-ribulokinase (EC 2.7.1.16); xylulokinase (EC 2.7.1.17); erythritol kinase (EC 2.7.1.27); glycerol kinase (EC 2.7.1.30); pantothenate kinase (EC 2.7.1.33); D-ribulokinase (EC 2.7.1.47); L-fucolokinase (EC 2.7.1.51); L-xylulokinase (EC 2.7.1.53); allose kinase (EC 2.7.1.55); 2-dehydro-3-deoxygalactonokinase (EC 2.7.1.58); N-acetylglucosamine kinase (EC 2.7.1.59); N-acylmannosamine kinase (EC 2.7.1.60); polyphosphate-glucose phosphotransferase (EC 2.7.1.63); beta-glucoside kinase (EC 2.7.1.85); acetate kinase (EC 2.7.2.1); butyrate kinase (EC 2.7.2.7); branched-chain-fatty-acid kinase (EC 2.7.2.14); and propionate kinase (EC 2.7.2.15).\nGroup V kinases (1,119 sequences) incorporate TIM β-barrel kinases. These enzymes include pyruvate kinase (EC 2.7.1.40).\nGroup VI kinases (885 sequences) incorporate GHMP kinases. These enzymes include galactokinase (EC 2.7.1.6); mevalonate kinase (EC 2.7.1.36); homoserine kinase (EC 2.7.1.39); L-arabinokinase (EC 2.7.1.46); fucokinase (EC 2.7.1.52); shikimate kinase (EC 2.7.1.71); 4-(cytidine 5′-diphospho)-2-C-methyl-D-erythriol kinase (EC 2.7.1.148); and phosphomevalonate kinase (EC 2.7.4.2).\nGroup VII kinases (1,843 sequences) incorporate AIR synthetase-like kinases. These enzymes include thiamine-phosphate kinase (EC 2.7.4.16) and selenide, water dikinase (EC 2.7.9.3).\nGroup VIII kinases (565 sequences) incorporate riboflavin kinases (565 sequences). These enzymes include riboflavin kinase (EC 2.7.1.26).\nGroup IX kinases (197 sequences) incorporate dihydroxyacetone kinases. These enzymes include glycerone kinase (EC 2.7.1.29).\nGroup X kinases (148 sequences) incorporate putative glycerate kinases. These enzymes include glycerate kinase (EC 2.7.1.31).\nGroup XI kinases (446 sequences) incorporate polyphosphate kinases. These enzymes include polyphosphate kinases (EC 2.7.4.1).\nGroup XII kinases (263 sequences) incorporate integral membrane kinases. Group XII comprises the dolichol kinase family. These enzymes include dolichol kinases (EC 2.7.1.108). Group XII further comprises the undecaprenol kinase family. These enzymes include undecaprenol kinases (EC 2.7.1.66).\nKinases play indispensable roles in numerous cellular metabolic and signaling pathways, and are among the best-studied enzymes at the structural, biochemical, and cellular level. Despite the fact that all kinases use the same phosphate donor (in most cases, ATP) and catalyze apparently the same phosphoryl transfer reaction, they display remarkable diversity in their structural folds and substrate recognition mechanisms. This probably is due largely to the diverse nature of the structures and properties of their substrates.\nMitogen-Activated Protein Kinase (MAPK)-Activated Protein Kinases (MK2 and MK3)\nDifferent groups of MAPK-activated protein kinases (MAP-KAPKs) have been defined downstream of mitogen-activated protein kinases (MAPKs). These enzymes transduce signals to target proteins that are not direct substrates of the MAPKs and, therefore, serve to relay phosphorylation-dependent signaling with MAPK cascades to diverse cellular functions. One of these groups is formed by the three MAPKAPKs: MK2, MK3 (also known as 3pK), and MK5 (also designated PRAK). Mitogen-activated protein kinase-activated protein kinase 2 (also referred to as “MAPKAPK2”, “MAPKAP-K2”, “MK2”) is a kinase of the serine/threonine (Ser/Thr) protein kinase family. MK2 is highly homologous to MK3 (approximately 75% amino acid identity). The kinase domains of MK2 and MK3 are most similar (approximately 35% to 40% identity) to calcium/calmodulin-dependent protein kinase (CaMK), phosphorylase b kinase, and the C-terminal kinase domain (CTKD) of the ribosomal S6 kinase (RSK) isoforms. The MK2 gene encodes two alternatively spliced transcripts of 370 amino acids (MK2A) and 400 amino acids (MK2B). The MK3 gene encodes one transcript of 382 amino acids. The MK2- and MK3 proteins are highly homologous, yet MK2A possesses a shorter C-terminal region. The C-terminus of MK2B contains a functional bipartite nuclear localization sequence (NLS) (Lys-Lys-Xaa-Xaa-Xaa-Xaa-Xaa-Xaa-Xaa-Xaa-Xaa-Xaa-Lys-Arg-Arg-Lys-Lys; SEQ ID NO: 21) that is not present in the shorter MK2A isoform, indicating that alternative splicing determines the cellular localization of the MK2 isoforms. MK3 possesses a similar nuclear localization sequence. The nuclear localization sequence found in both MK2B and MK3 encompasses a D domain (Leu-Leu-Lys-Arg-Arg-Lys-Lys; SEQ ID NO: 22), which was shown to mediate the specific interaction of MK2B and MK3 with p38α and p38β. MK2B and MK3 also possess a functional nuclear export signal (NES) located N-terminal to the NLS and D domain. The NES in MK2B is sufficient to trigger nuclear export following stimulation, a process which may be inhibited by leptomycin B. The sequence N-terminal to the catalytic domain in MK2 and MK3 is proline rich and contains one (MK3) or two (MK2) putative Src homology 3 (SH3) domain-binding sites, which studies have shown, for MK2, to mediate binding to the SH3 domain of c-Abl in vitro. Recent studies suggest that this domain is involved in MK2-mediated cell migration.\nMK2B and MK3 are located predominantly in the nucleus of quiescent cells while MK2A is present in the cytoplasm. Both MK2B and MK3 are rapidly exported to the cytoplasm via a chromosome region maintenance protein (CRM1)-dependent mechanism upon stress stimulation. Nuclear export of MK2B appears to be mediated by kinase activation, as phosphomimetic mutation of Thr334 within the activation loop of the kinase enhances the cytoplasmic localization of MK2B. Without being limited by theory, it is thought that MK2B and MK3 may contain a constitutively active nuclear localization signal (NLS) and a phosphorylation-regulated nuclear export signal (NES).\nMK2 and MK3 appear to be expressed ubiquitously, with increased relative expression in the heart, lungs, kidney, reproductive organs (mammary and testis), skin and skeletal muscle tissues, as well as in immune-related cells such as white blood cells/leukocytes and dendritic cells.\nActivation of MK2 and MK3 Activity\nVarious activators of p38α and p38β potently stimulate MK2 and MK3 activity. p38 mediates the in vitro and in vivo phosphorylation of MK2 on four proline-directed sites: Thr25, Thr222, Ser272, and Thr334. Of these sites, only Thr25 is not conserved in MK3. Without being limited by theory, while the function of phosphorylated Thr25 is unknown, its location between the two SH3 domain-binding sites suggests that it may regulate protein-protein interactions. Thr222 in MK2 (Thr201 in MK3) is located in the activation loop of the kinase domain and has been shown to be essential for MK2 and MK3 kinase activity. Thr334 in MK2 (Thr313 in MK3) is located C-terminal to the catalytic domain and is essential for kinase activity. The crystal structure of MK2 has been resolved and, without being limited by theory, suggests that Thr334 phosphorylation may serve as a switch for MK2 nuclear import and export. Phosphorylation of Thr334 also may weaken or interrupt binding of the C terminus of MK2 to the catalytic domain, exposing the NES and promoting nuclear export.\nStudies have shown that while p38 is capable of activating MK2 and MK3 in the nucleus, experimental evidence suggests that activation and nuclear export of MK2 and MK3 are coupled by a phosphorylation-dependent conformational switch that also dictates p38 stabilization and localization, and the cellular location of p38 itself is controlled by MK2 and possibly MK3. Additional studies have shown that nuclear p38 is exported to the cytoplasm in a complex with MK2 following phosphorylation and activation of MK2. The interaction between p38 and MK2 may be important for p38 stabilization since studies indicate that p38 levels are low in MK2-deficient cells and expression of a catalytically inactive MK2 protein restores p38 levels.\nSubstrates and Functions\nMK2 shares many substrates with MK3. Both enzymes have comparable substrate preferences and phosphorylate peptide substrates with similar kinetic constants. The minimum sequence required for efficient phosphorylation by MK2 was found to be Hyd-Xaa-Arg-Xaa-Xaa-pSer/pThr (SEQ ID NO: 22), where Hyd is a bulky, hydrophobic residue.\nAccumulating evidence has shown that MK2 phophorylates a variety of proteins, which include, but are not limited to, 5-Lipooxygenase (ALOX5), Cell Division Cycle 25 Homolog B (CDC25B), Cell Division Cycle 25 Homolog C (CDC25C), Embryonic Lethal, Abnormal Vision, Drosophila-Like 1 (ELAVL1), Heterogeneous Nuclear Ribonucleoprotein A0 (HNRNPA0), Heat Shock Factor protein 1 (HSF1), Heat Shock Protein Beta-1 (HSPB1), Keratin 18 (KRT18), Keratin 20 (KRT20), LIM domain kinase 1 (LIMK1), Lymphocyte-specific protein 1 (LSP1), Polyadenylate-Binding Protein 1 (PABPC1), Poly(A)-specific Ribonuclease (PARN), CAMP-specific 3′,5′-cyclic Phosphodiesterase 4A (PDE4A), RCSD domain containing 1 (RCSD1), Ribosomal protein S6 kinase, 90 kDa, polypeptide 3 (RPS6KA3), TGF-beta activated kinase 1/MAP3K7 binding protein 3 (TAB3), and Tristetraprolin (TTP/ZFP36).\nHeat-Shock Protein Beta-1 (also termed HSPB1 or HSP27) is a stress-inducible cytosolic protein that is ubiquitously present in normal cells and is a member of the small heat-shock protein family. The synthesis of HSPB1 is induced by heat shock and other environmental or pathophysiologic stresses, such as UV radiation, hypoxia and ischemia. Besides its putative role in thermoresistance, HSPB1 is involved in the survival and recovery of cells exposed to stressful conditions.\nExperimental evidence supports a role for p38 in the regulation of cytokine biosynthesis and cell migration. The targeted deletion of the mk2 gene in mice suggested that although p38 mediates the activation of many similar kinases, MK2 seems to be the key kinase responsible for these p38-dependent biological processes. Loss of MK2 leads (i) to a defect in lipopolysaccharide (LPS)-induced synthesis of cytokines such as tumor necrosis factor alpha (TNF-α), interleukin-6 (IL-6), and gamma interferon (IFN-γ) and (ii) to changes in the migration of mouse embryonic fibroblasts, smooth muscle cells, and neutrophils.\nConsistent with a role for MK2 in inflammatory and immune responses, MK2-deficient mice showed increased susceptibility to Listeria monocytogenes infection and reduced inflammation-mediated neuronal death following focal ischemia. Since the levels of p38 protein also are reduced significantly in MK2-deficient cells, it was necessary to distinguish whether these phenotypes were due solely to the loss of MK2. To achieve this, MK2 mutants were expressed in MK2-deficient cells, and the results indicated that the catalytic activity of MK2 was not necessary to restore p38 levels but was required to regulate cytokine biosynthesis.\nKnockout or knockdown studies of MK2 provide strong support that activated MK2 enhances stability of IL-6 mRNA through phosphorylation of proteins interacting with the AU-rich 3′ untranslated region of IL-6 mRNA. In particular, it has been shown that MK2 is principally responsible for phosphorylation of hnRNPA0, an mRNA-binding protein that stabilizes IL-6 RNA. In addition, several additional studies investigating diverse inflammatory diseases have found that levels of pro-inflammatory cytokines, such as IL-6, IL-1β, TNF-α and IL-8, are increased in induced sputum from patients with stable chronic obstructive pulmonary disease (COPD) or from the alveolar macrophages of cigarette smokers (Keatings V. et al, Am J Resp Crit Care Med, 1996, 153:530-534; Lim, S. et al., J Respir Crit Care Med, 2000, 162:1355-1360).\nRegulation of mRNA Translation.\nPrevious studies using MK2 knockout mice or MK2-deficient cells have shown that MK2 increases the production of inflammatory cytokines, including TNF-α, IL-1, and IL-6, by increasing the rate of translation of its mRNA. No significant reductions in the transcription, processing, and shedding of TNF-α could be detected in MK2-deficient mice. The p38 pathway is known to play an important role in regulating mRNA stability, and MK2 represents a likely target by which p38 mediates this function. Studies utilizing MK2-deficient mice indicated that the catalytic activity of MK2 is necessary for its effects on cytokine production and migration, suggesting that, without being limited by theory, MK2 phosphorylates targets involved in mRNA stability. Consistent with this, MK2 has been shown to bind and/or phosphorylate the heterogeneous nuclear ribonucleoprotein (hnRNP) A0, tristetraprolin (TTP), the poly(A)-binding protein PABP1, and HuR, a ubiquitously expressed member of the ELAV (Embryonic-Lethal Abnormal Visual in Drosophila melanogaster) family of RNA-binding protein. These substrates are known to bind or copurify with mRNAs that contain AU-rich elements in the 3′ untranslated region, suggesting that MK2 may regulate the stability of AU-rich mRNAs such as TNF-α. It currently is unknown whether MK3 plays a similar role, but LPS treatment of MK2-deficient fibroblasts completely abolished hnRNP A0 phosphorylation, suggesting that MK3 is not able to compensate for the loss of MK2.\nMK3 participates with MK2 in phosphorylation of the eukaryotic elongation factor 2 (eEF2) kinase. eEF2 kinase phosphorylates and inactivates eEF2. eEF2 activity is critical for the elongation of mRNA during translation, and phosphorylation of eEF2 on Thr56 results in the termination of mRNA translation. MK2 and MK3 phosphorylation of eEF2 kinase on Ser377 suggests that these enzymes may modulate eEF2 kinase activity and thereby regulate mRNA translation elongation.\nTranscriptional Regulation by MK2 and MK3\nNuclear MK2, similar to many MKs, contributes to the phosphorylation of cAMP response element binding (CREB), Activating Transcription Factor-1 (ATF-1), serum response factor (SRF), and transcription factor ER81. Comparison of wild-type and MK2-deficient cells revealed that MK2 is the major SRF kinase induced by stress, suggesting a role for MK2 in the stress-mediated immediate-early response. Both MK2 and MK3 interact with basic helix-loop-helix transcription factor E47 in vivo and phosphorylate E47 in vitro. MK2-mediated phosphorylation of E47 was found to repress the transcriptional activity of E47 and thereby inhibit E47-dependent gene expression, suggesting that MK2 and MK3 may regulate tissue-specific gene expression and cell differentiation.\nOther Targets of MK2 and MK3\nSeveral other MK2 and MK3 substrates also have been identified, reflective of the diverse functions of MK2 and MK3 in several biological processes. The scaffolding protein 14-3-3ζ is a physiological MK2 substrate. Studies indicate that 14-3-3ζ interacts with a number of components of cell signaling pathways, including protein kinases, phosphatases, and transcription factors. Additional studies have shown that MK2-mediated phosphorylation of 14-3-3ζ on Ser58 compromises its binding activity, suggesting that MK2 may affect the regulation of several signaling molecules normally regulated by 14-3-3ζ.\nAdditional studies have shown that MK2 also interacts with and phosphorylates the p16 subunit of the seven-member Arp2 and Arp3 complex (p16-Arc) on Ser77. p16-Arc has roles in regulating the actin cytoskeleton, suggesting that MK2 may be involved in this process. Further studies have shown that the small heat shock protein HSPB1, lymphocyte-specific protein LSP-1, and vimentin are phosphorylated by MK2. HSPB1 forms large oligomers which may act as molecular chaperones and protect cells from heat shock and oxidative stress. Upon phosphorylation, HSPB1 loses its ability to form large oligomers and is unable to block actin polymerization, suggesting that MK2-mediated phosphorylation of HSPB1 serves a homeostatic function aimed at regulating actin dynamics that otherwise would be destabilized during stress. MK3 also was shown to phosphorylate HSPB1 in vitro and in vivo, but its role during stressful conditions has not yet been elucidated.\nIt was also shown that HSPB1 binds to polyubiquitin chains and to the 26S proteasome in vitro and in vivo. The ubiquitin-proteasome pathway is involved in the activation of transcription factor NF-kappa B (NF-κB) by degrading its main inhibitor, I kappa B-alpha (IκB-alpha), and it was shown that overexpresion of HSPB1 increases NF-kappaB (NF-κB) nuclear relocalization, DNA binding, and transcriptional activity induced by etoposide, TNF-alpha, and Interleukin-1 beta (IL-1β). Additionally, previous studies have suggested that HSPB1, under stress conditions, favors the degradation of ubiquitinated proteins, such as phosphorylated I kappa B-alpha (IκB-alpha); and that this function of HSPB1 accounts for its anti-apoptotic properties through the enhancement of NF-kappa B (NF-κB) activity (Parcellier, A. et al., Mol Cell Biol, 23(16): 5790-5802, 2003).\nMK2 and MK3 also may phosphorylate 5-lipoxygenase. 5-lipoxygenase catalyzes the initial steps in the formation of the inflammatory mediators, leukotrienes. Tyrosine hydroxylase, glycogen synthase, and Akt also were shown to be phosphorylated by MK2.\nMK2 phosphorylates the tumor suppressor protein tuberin on Ser1210, creating a docking site for 14-3-3ζ. Tuberin and hamartin normally form a functional complex that negatively regulates cell growth by antagonizing mTOR-dependent signaling, suggesting that p38-mediated activation of MK2 may regulate cell growth by increasing 14-3-3ζ binding to tuberin.\nAccumulating evidence has suggested that the reciprocal crosstalk between the p38 MAPK-pathway and signal transducer and activator of transcription 3 (STAT3)-mediated signal-transduction forms a critical axis successively activated in lipopolysaccharide (LPS) challenge models. It was shown that the balanced activation of this axis is essential for both induction and propagation of the inflammatory macrophage response as well as for the control of the resolution phase, which is largely driven by IL-10 and sustained STAT3 activation (Bode, J. et al., Cellular Signalling, 24: 1185-1194, 2012). In addition, another study has shown that MK2 controls LPS-inducible IFNβ gene expression and subsequent IFNβ-mediated activation of STAT3 by neutralizing negative regulatory effects of MK3 on LPS-induced p65 and IRF3-mediated signaling. In mk2/3 knockout macrophages, IFNβ-dependent STAT3 activation occurred independently from IL-10, because, in contrast to IFNβ-, impaired IL-10 expression is not restored upon additional deletion of MK3 in mk2/3 knockout macrophages (Ehlting, C. et al., J. Biol. Chem., 285(27): 24113-24124).\nKinase Inhibition\nThe eukaryotic protein kinases constitute one of the largest superfamilies of homologous proteins that are related by virtue of their catalytic domains. Most related protein kinases are specific for either serine/threonine or tyrosine phosphorylation. Protein kinases play an integral role in the cellular response to extracellular stimuli. Thus, stimulation of protein kinases is considered to be one of the most common activation mechanisms in signal transduction systems. Many substrates are known to undergo phosphorylation by multiple protein kinases, and a considerable amount of information on primary sequence of the catalytic domains of various protein kinases has been published. These sequences share a large number of residues involved in ATP binding, catalysis, and maintenance of structural integrity. Most protein kinases possess a well conserved 30-32 kDa catalytic domain.\nStudies have attempted to identify and utilize regulatory elements of protein kinases. These regulatory elements include inhibitors, antibodies, and blocking peptides.\nInhibitors\nEnzyme inhibitors are molecules that bind to enzymes thereby decreasing enzyme activity. The binding of an inhibitor may stop a substrate from entering the active site of the enzyme and/or hinder the enzyme from catalyzing its reaction (as in inhibitors directed at the ATP biding site of the kinase). Inhibitor binding is either reversible or irreversible. Irreversible inhibitors usually react with the enzyme and change it chemically (e.g., by modifying key amino acid residues needed for enzymatic activity) so that it no longer is capable of catalyzing its reaction. In contrast, reversible inhibitors bind non-covalently, and different types of inhibition are produced depending on whether these inhibitors bind the enzyme, the enzyme-substrate complex, or both.\nEnzyme inhibitors often are evaluated by their specificity and potency. The term “specificity” as used in this context refers to the selective attachment of an inhibitor or its lack of binding to other proteins. The term “potency” as used herein refers to an inhibitor's dissociation constant, which indicates the concentration of inhibitor needed to inhibit an enzyme.\nInhibitors of protein kinases have been studied for use as a tool in protein kinase activity regulation. Inhibitors have been studied for use with, for example, cyclin-dependent (Cdk) kinase, MAP kinase, serine/threonine kinase, Src Family protein tyrosine kinase, tyrosine kinase, calmodulin (CaM) kinase, casein kinase, checkpoint kinase (Chk1), glycogen synthase kinase 3 (GSK-3), c-Jun N-terminal kinase (JNK), mitogen-activated protein kinase 1 (MEK), myosin light chain kinase (MLCK), protein kinase A, Akt (protein kinase B), protein kinase C, protein kinase G, protein tyrosine kinase, Raf kinase, and Rho kinase.\nSmall-Molecule MK2 Inhibitors\nWhile individual inhibitors that target MK2 with at least modest selectivity with respect to other kinases have been designed, it has been difficult to create compounds with favorable solubility and permeability. As a result, there are relatively few biochemically efficient MK2 inhibitors that have advanced to in vivo pre-clinical studies (Edmunds, J. and Talanian, MAPKAP Kinase 2 (MK2) as a Target for Anti-inflammatory Drug Discovery. In Levin, J and Laufer, S (Ed.), RSC Drug Discovery Series No. 26, p 158-175, the Royal Society of Chemistry, 2012; incorporated by reference in its entirety).\nThe majority of disclosed MK2 inhibitors are classical type I inhibitors as revealed by crystallographic or biochemical studies. As such, they bind to the ATP site of the kinase and thus compete with intra-cellular ATP (estimated concentration 1 mM-5 mM) to inhibit phosphorylation and activation of the kinase. Representative examples of such small-molecule MK2 inhibitors include, but are not limited to,\n\nBlocking Peptides\nA peptide is a chemical compound composed of a chain of two or more amino acids whereby the carboxyl group of one amino acid in the chain is linked to the amino group of the other via a peptide bond. Peptides have been used inter alia in the study of protein structure and function. Synthetic peptides may be used inter alia as probes to see where protein-peptide interactions occur. Inhibitory peptides may be used inter alia in clinical research to examine the effects of peptides on the inhibition of protein kinases, cancer proteins and other disorders.\nThe use of several blocking peptides has been studied. For example, extracellular signal-regulated kinase (ERK), a MAPK protein kinase, is essential for cellular proliferation and differentiation. The activation of MAPKs requires a cascade mechanism whereby MAPK is phosphorylated by an upstream MAPKK (MEK) which then, in turn, is phosphorylated by a third kinase MAPKKK (MEKK). The ERK inhibitory peptide functions as a MEK decoy by binding to ERK.\nThe synthetic peptide autocamtide-2 related inhibitory peptide (AIP), is a highly specific and potent inhibitor of Ca2+/calmodulin-dependent protein kinase II (CaMKII). AIP is a non-phosphorylatable analog of autocamtide-2, a highly selective peptide substrate for CaMKII. AIP inhibits CaMKII with an IC50 of 100 nM (IC50 is the concentration of an inhibitor required to obtain 50% inhibition). The AIP inhibition is non-competitive with respect to syntide-2 (CaMKII peptide substrate) and ATP but competitive with respect to autocamtide-2. The inhibition is unaffected by the presence or absence of Ca2+/calmodulin. CaMKII activity is inhibited completely by AIP (1 μM) while PKA, PKC and CaMKIV are not affected.\nCell division protein kinase 5 (Cdk5) inhibitory peptide (CIP) is another blocking peptide. Cdk5 phosphorylates the microtubule protein tau at Alzheimer's Disease-specific phospho-epitopes when it associates with p25. p25 is a truncated activator, which is produced from the physiological Cdk5 activator p35 upon exposure to amyloid β peptides. Upon neuronal infections with CIP, CIPs selectively inhibit p25/Cdk5 activity and suppress the aberrant tau phosphorylation in cortical neurons. The reasons for the specificity demonstrated by CIP are not fully understood.\nAdditional blocking peptides have been studied for inhibiting extracellular-regulated kinase 2 (ERK2), ERK3, p38/HOG1, protein kinase C, casein kinase II, Ca2+/calmodulin kinase IV, casein kinase II, Cdk4, Cdk5, DNA-dependent protein kinase (DNA-PK), serine/threonine-protein kinase PAK3, phosphoinositide (PI)-3 kinase, PI-5 kinase, PSTAIRE (the cdk highly conserved sequence), ribosomal S6 kinase, GSK-4, germinal center kinase (GCK), SAPK (stress-activated protein kinase), SEK1 (stress signaling kinase), and focal adhesion kinase (FAK).\nProtein Substrate-Competitive Inhibitors\nMost of the protein kinase inhibitors developed to date are ATP competitors. This type of molecule competes for the ATP binding site of the kinase and often shows off-target effects due to serious limitations in its specificity. The low specificity of these inhibitors is due to the fact that the ATP binding site is highly conserved among diverse protein kinases. Non-ATP competitive inhibitors, on the other hand, such as substrate competitive inhibitors, are expected to be more specific as the substrate binding sites have a certain degree of variability among the various protein kinases.\nAlthough substrate competitive inhibitors usually have a weak binding interaction with the target enzyme in vitro, studies have shown that chemical modifications can improve the specific biding affinity and the in vivo efficacy of substrate inhibitors (Eldar-Finkelman, H. et al., Biochim, Biophys. Acta, 1804(3):598-603, 2010). In addition, substrate competitive inhibitors show better efficacy in cells than in cell-free conditions in many cases (van Es, J. et al., Curr. Opin. Gent. Dev. 13:28-33, 2003).\nIn an effort to enhance specificity and potency in protein kinase inhibition, bisubstrate inhibitors also have been developed. Bisubstrate inhibitors, which consist of two conjugated fragments, each targeted to a different binding site of a bisubstrate enzyme, form a special group of protein kinase inhibitors that mimic two natural substrates/ligands and that simultaneously associate with two regions of given kinases. The principle advantage of bisubstrate inhibitors is their ability to generate more interactions with the target enzyme that could result in improved affinity and selectivity of the conjugates, when compared with single-site inhibitors. Examples of bisubstrate inhibitors include, but are not limited to, nucleotide-peptide conjugates, adenosine derivative-peptide conjugates, and conjugates of peptides with potent ATP-competitive inhibitors.\nProtein Transduction Domains (PTD)/Cell Permeable Proteins (CPP)\nThe plasma membrane presents a formidable barrier to the introduction of macromolecules into cells. For nearly all therapeutics to exert their effects, at least one cellular membrane must be traversed. Traditional small molecule pharmaceutical development relies on the chance discovery of membrane permeable molecules with the ability to modulate protein function. Although small molecules remain the dominant therapeutic paradigm, many of these molecules suffer from lack of specificity, side effects, and toxicity. Information-rich macromolecules, which have protein modulatory functions far superior to those of small molecules, can be created using rational drug design based on molecular, cellular, and structural data. However, the plasma membrane is impermeable to most molecules of size greater than 500 Da. Therefore, the ability of cell penetrating peptides, such as the basic domain of Trans-Activator of Transcription (Tat), to cross the cell membrane and deliver macromolecular cargo in vivo, can greatly facilitate the rational design of therapeutic proteins, peptides, and nucleic acids.\nProtein transduction domains (PTDs) are a class of peptides capable of penetrating the plasma membrane of mammalian cells and of transporting compounds of many types and molecular weights across the membrane. These compounds include effector molecules, such as proteins, DNA, conjugated peptides, oligonucleotides, and small particles such as liposomes. When PTDs are chemically linked or fused to other proteins, the resulting fusion peptides still are able to enter cells. Although the exact mechanism of transduction is unknown, internalization of these proteins is not believed to be receptor-mediated or transporter-mediated. PTDs are generally 10-16 amino acids in length and may be grouped according to their composition, such as, for example, peptides rich in arginine and/or lysine.\nThe use of PTDs capable of transporting effector molecules into cells has become increasingly attractive in the design of drugs as they promote the cellular uptake of cargo molecules. These cell-penetrating peptides, generally categorized as amphipathic (meaning having both a polar and a nonpolar end) or cationic (meaning of or relating to containing net positively charged atoms) depending on their sequence, provide a non-invasive delivery technology for macromolecules. PTDs often are referred to as “Trojan peptides”, “membrane translocating sequences”, or “cell permeable proteins” (CPPs). PTDs also may be used to assist novel HSPB1 kinase inhibitors to penetrate cell membranes. (see, e.g., U.S. Pat. No. 8,536,303, and U.S. Pat. No. 8,741,849, the contents of each of which are incorporated by reference in their entirety herein).\nViral PTD Containing Proteins\nThe first proteins to be described as having transduction properties were of viral origin. These proteins still are the most commonly accepted models for PTD action. The HIV-1 Transactivator of Transcription (Tat) and HSV-1 VP 22 protein are the best characterized viral PTD containing proteins.\nTat (HIV-1 trans-activator gene product) is an 86-amino acid polypeptide, which acts as a powerful transcription factor of the integrated HIV-1 genome. Tat acts on the viral genome, stimulating viral replication in latently infected cells. The translocation properties of the Tat protein enable it to activate quiescent infected cells, and it may be involved in priming of uninfected cells for subsequent infection by regulating many cellular genes, including cytokines. The minimal PTD of Tat is the 9 amino acid protein sequence RKKRRQRRR (TAT49-57; SEQ ID NO: 20). Studies utilizing a longer fragment of Tat demonstrated successful transduction of fusion proteins up to 120 kDa. The addition of multiple Tat-PTDs as well as synthetic Tat derivatives has been demonstrated to mediate membrane translocation. Tat PTD containing fusion proteins have been used as therapeutic moieties in experiments involving cancer, transporting a death-protein into cells, and disease models of neurodegenerative disorders.\nThe mechanism used by transducing peptides to permeate cell membranes has been the subject of considerable interest in recent years, as researchers have sought to understand the biology behind transduction. Early reports that Tat transduction occurred by a nonendocytic mechanism have largely been dismissed as artifactual though other cell-penetrating peptides might be taken up by way of direct membrane disruption. The recent findings that transduction of Tat and other PTDs occurs by way of macropinocytosis, a specialized form of endocytosis, has created a new paradigm in the study of these peptides. Enhanced knowledge of the mechanism of transduction helped improve transduction efficiency with the ultimate goal of clinical success (Snyder E. and Dowdy, S., Pharm Res., 21(3):389-393, 2004).\nThe current model for Tat-mediated protein transduction is a multistep process that involves binding of Tat to the cell surface, stimulation of macropinocytosis, uptake of Tat and cargo into macropinosomes, and endosomal escape into the cytoplasm. The first step, binding to the cell surface, is thought to be through ubiquitous glycan chains on the cell surface. Stimulation of macropinocytosis by Tat occurs by an unknown mechanism that might include binding to a cell surface protein or might occur by way of proteoglycans or glycolipids. Uptake by way of macropinocytosis, a form of fluid phase endocytosis used by all cell types, is required for Tat and polyarginine transduction. The final step in Tat transduction is escape from macropinosomes into the cytoplasm; this process is likely to be dependent on the pH drop in endosomes that, along with other factors, facilitates a perturbation of the membrane by Tat and release of Tat and its cargo (i.e. peptide, protein or drug etc.) to the cytoplasm (Snyder E. and Dowdy, S., Pharm Res., 21(3):389-393, 2004).\nVP22 is the HSV-1 tegument protein, a structural part of the HSV virion. VP22 is capable of receptor independent translocation and accumulates in the nucleus. This property of VP22 classifies the protein as a PTD containing peptide. Fusion proteins comprising full length VP22 have been translocated efficiently across the plasma membrane.\nHomeoproteins with Intercellular Translocation Properties\nHomeoproteins are highly conserved, transactivating transcription factors involved in morphological processes. They bind to DNA through a specific sequence of 60 amino acids. The DNA-binding homeodomain is the most highly conserved sequence of the homeoprotein. Several homeoproteins have been described as exhibiting PTD-like activity; they are capable of efficient translocation across cell membranes in an energy-independent and endocytosis-independent manner without cell type specificity.\nThe Antennapedia protein (Antp) is a trans-activating factor capable of translocation across cell membranes; the minimal sequence capable of translocation is a 16 amino acid peptide corresponding to the third helix of the protein's homeodomain (HD). The internalization of this helix occurs at 4° C., suggesting that this process is not endocytosis dependent. Peptides of up to 100 amino acids produced as fusion proteins with AntpHD penetrate cell membranes.\nOther homeodomains capable of translocation include Fushi tarazu (Ftz) and Engrailed (En) homeodomain. Many homeodomains share a highly conserved third helix.\nHuman PTDs\nHuman PTDs may circumvent potential immunogenicity issues upon introduction into a human patient. Peptides with PTD sequences include: Hoxa-5, Hox-A4, Hox-B5, Hox-B6, Hox-B7, HOX-D3, GAX, MOX-2, and FtzPTD. These proteins all share the sequence found in AntpPTD. Other PTDs include Islet-1, Interleukin-1 (IL-1), Tumor Necrosis Factor (TNF), and the hydrophobic sequence from Kaposi-fibroblast growth factor or Fibroblast Growth Factor-4 (FGF-4) signal peptide, which is capable of energy-, receptor-, and endocytosis-independent translocation. Unconfirmed PTDs include members of the Fibroblast Growth Factor (FGF) family. FGFs are polypeptide growth factors that regulate proliferation and differentiation of a wide variety of cells. Several publications have reported that basic fibroblast growth factor (FGF-2) exhibits an unconventional internalization similar to that of VP-22, Tat, and homeodomains. It has also been reported that acidic FGF (FGF-1) translocated cell membranes at temperatures as low as 4° C. However, no conclusive evidence exists about the domain responsible for internalization or the translocation properties of fusion proteins (Beerens, A. et al., Curr Gene Ther., 3(5):486-494, 2003).\nSynthetic PTDs\nSeveral peptides have been synthesized in an attempt to create more potent PTDs and to elucidate the mechanisms by which PTDs transport proteins across cell membranes. Many of these synthetic PTDs are based on existing and well documented peptides, while others are selected for their basic residues and/or positive charges, which are thought to be crucial for PTD function. A few of these synthetic PTDs showed better translocation properties than the existing ones (Beerens, A. et al., Curr Gene Ther., 3(5):486-494, 2003). Exemplary Tat-derived synthetic PTDs include, for example, but are not limited to, WLRRIKAWLRRIKA (SEQ ID NO: 12); WLRRIKA (SEQ ID NO: 13); YGRKKRRQRRR (SEQ ID NO: 14); WLRRIKAWLRRI (SEQ ID NO: 15); FAKLAARLYR (SEQ ID NO: 16); KAFAKLAARLYR (SEQ ID NO: 17); and HRRIKAWLKKI (SEQ ID NO: 18).\nCompositions Comprising PTDs Fused to MK2 Inhibitor Peptide Therapeutic Domains (TD)\nSeveral MK2 inhibitor peptides (TD) have been synthesized, fused to synthetic PTDs and the use of compositions comprising these fused polypeptides has been studied. These polypeptides include, but are not limited to, YARAAARQARAKALARQLGVAA (SEQ ID NO: 1; MMI-0100), YARAAARQARAKALNRQLGVA (SEQ ID NO: 19; MMI-0200), FAKLAARLYRKALARQLGVAA (SEQ ID NO: 3; MMI-0300), KAFAKLAARLYRKALARQLGVAA (SEQ ID NO: 4; MMI-0400), HRRIKAWLKKIKALARQLGVAA (SEQ ID NO: 7; MMI-0500), YARAAARDARAKALNRQLAVAA (SEQ ID NO: 23; MMI-0600), and YARAAARQARAKALNRQLAVA (SEQ ID NO: 24; MMI-0600-2). Both in vitro and in vivo studies have shown that these polypeptides can be useful in the treatment of various diseases, disorders and conditions. These include, without limitation, hyperplasia and neoplasm (U.S. Pat. Nos. 8,536,303 and 8,741,849) inflammatory disorders (U.S. application Ser. No. 12/634,476 and U.S. application Ser. No. 13/934,933), adhesions (U.S. application Ser. No. 12/582,516), failure of a vascular graft due to neospasm (U.S. application Ser. No. 13/114,872), improving neurite outgrowth (U.S. application Ser. No. 12/844,815), cutaneous scarring (U.S. application Ser. No. 13/829,876), coronary artery bypass vascular graft failure (U.S. application Ser. No. 13/700,087) and interstitial lung disease and pulmonary fibrosis (U.S. application Ser. No. 13/445,759).\nPeptide compositions present a number of particular challenges to formulation scientists (R. W. Payne and M. C. Manning, “Peptide formulation: challenges and strategies,” Innovations in Pharmaceutical Technology, 64-68 (2009)). First, since peptides do not have a globular structure that can sequester reactive groups, the side chains of nearly all residues in a peptide are fully solvent exposed, and can exhibit chemical degradation through hydrolytic reactions, for example, oxidation and deamidation. Second, the conformation in aqueous solution may have little similarity to the structure found when bound to a receptor. Third, many peptides tend to be monomeric at very low concentration, but may self-assemble as the concentration is increased and behave as if in a highly associated state, but these structures are too transient or fluxional to provide any increase in long-term stability. Fourth, the propensity of peptides to self-associate is connected with their physical instability, meaning their likelihood of forming aggregates. Moreover, excipients present in a peptide formulation can chemically degrade, interact with various surfaces during manufacturing, interact with the container or closure, or interact with the peptide itself, thereby negatively affecting critical properties of the preparation (Lars Hovgaard, and Sven Frokjaer, “Pharmaceutical Formulation Development of Peptides and Proteins, 2nd Ed., CRC Press (2012) pp. 212-213).\nTrkB is Highly Expressed in NSCLC\nAberrant regulation in the invasion of cancer cells is closely associated with their metastatic potential. (Zhang, S. et al., “TrkB is highly expressed in NSCLC and mediates BDNF-induced activation of Pyk2 signaling and the invasion of A549 cells,” BMC Cancer 2010, 10: 43). Tropomysin-related kinase B (TrkB), a member of the Trk family, functions as a receptor tyrosine kinase and is considered to facilitate tumor metastasis. Id. Brain-derived neurotrophic factor (BDNF), the primary ligand binding to TrkB, results in the regulation of various cellular activities in neuroblastoma, such as cell differentiation, apoptosis, and invasion. Id. TrkB is up-regulated in a variety of primary human tumors, including neuroblastoma and ovarian cancer. Id. Enhanced TrkB signaling promotes cell survival in an anchorage dependent manner. Id. When activated by BDNF, TrkB leads to the activation of downstream signaling molecules, such as phosphoinositide-3 kinase/protein kinase B (PI3K/Akt), which induces the differential regulation of apoptosis and metastasis. Id.\nProline-rich tyrosine kinase 2 (Pyk2) is an extensively expressed non-receptor tyrosine kinase that integrates signals from receptor tyrosine kinases and intracellular signaling molecules in such essential cellular processes as cell differentiation, proliferation and migration. Id. Pyk2 is rapidly tyrosine phosphorylated in response to various extracellular signals, and activated Pyk2 signaling promotes cell survival and migration in an anchorage-independent manner. Id. The tyrosine 402 (Tyr402) of Pyk2 serves as the primary autophosphorylation site essential for Pyk2 activity and function. Id. High activity of Tyr402 is found in tumor cells with a more invasive and metastatic phenotype. Id.\nHigh expression of TrkB has been reported in NSCLC, particularly correlated with lymph node metastasis and TNM stage. Id. TrkB-siRNA interrupted BDNF-promoted Pyk2 and extracellular regulating kinase (ERK) activation and invasion of lung adenocarcinoma A549 cells. Id. Moreover, Pyk2-siRNA inhibited BDNF-associated ERK phosphorylation and cell invastion. Id. The investigators considered these data as evidence that TrkB/Pyk2/ERK signaling mediates BDNF-induced invasion of A549 cells, and suggested that suppression of TrkB may provide a helpful target for inhibitory therapies of metastasis in NSCLC. Id.\nThe described invention offers an approach to treat non-small cell lung carcinoma using cell-penetrating, peptide-based inhibitors of MK2, that may be effective to overcome tumor chemoresistance, to enhance tumor chemosensitivity and to slow progression of the non-small cell lung cancer solid tumor comprising the population of tumor cells."}
{"text": "The present invention relates generally to medical devices and more particularly to delivery systems for self-expandable medical devices.\nSelf-expanding medical devices are used by physicians to treat numerous conditions using minimally invasive procedures. Examples of self-expanding medical devices include stents, stent-grafts, filters, valves, etc. Typically, self-expanding medical devices are made from an elastic structure that may be compressed into a low profile state that can be passed through vessels in a patient with minimal trauma. Once at the desired treatment site, the self-expanding medical device is released and self-expands like a spring until it contacts a tissue wall which prevents further expansion. Common materials that are used in self-expanding medical devices include nitinol and stainless steel, although other materials are also possible.\nOne type of self-expanding medical device that has become especially common is intraluminal stents. Stents are used to treat various organs, such as the vascular system, colon, biliary tract, urinary tract, esophagus, trachea and the like. For example, stents are commonly used to treat blockages, occlusions, narrowing ailments and other similar problems that restrict flow through a passageway. One area where stents are commonly used for treatment involves implanting an endovascular stent into the vascular system in order to improve or maintain blood flow through narrowed arteries. However, stents are also used in other treatments as well, such as the treatment of aneurysms. Stents have been shown to be useful in treating various vessels throughout the vascular system, including both coronary vessels and peripheral vessels (e.g., carotid, brachial, renal, iliac and femoral). In addition, stents have been used in other body vessels as well, such as the digestive tract.\nThe use of stents in coronary and peripheral vessels has drawn particular attention from the medical community because of the growing number of people each year that suffer from vasculature problems associated with stenosis (i.e., narrowing of a vessel). This has led to an increased demand for medical procedures to treat such problems. The widespread frequency of heart problems and other vasculature problems may be due to a number of societal changes, including the tendency of people to exercise less and the prevalence of unhealthy diets, in conjunction with the fact that people generally have longer life spans now than previous generations. Stents have become a popular alternative for treating vascular stenosis because stenting procedures are considerably less invasive than conventional procedures. For example, stenosis of the coronary arteries was traditionally treated with bypass surgery. In general, bypass surgery involves splitting the chest bone to open the chest cavity and grafting a replacement vessel onto the heart to bypass the blocked, or stenosed, artery. However, coronary bypass surgery is a very invasive procedure that is risky and requires a long recovery time for the patient. Vascular stents are also being more widely used to treat many different peripheral arteries due to the minimally invasive nature of stenting procedures. To address the growing demand for minimally invasive medical procedures for the treatment of coronary arteries, peripheral arteries and other passageway problems, the medical community has begun to turn away from conventional invasive procedures like bypass surgery and increasingly the treatment of choice now involves a variety of stenting procedures.\nMany different types of stents and stenting procedures are possible. In general, however, stents are typically designed as tubular support structures that may be inserted percutaneously and transluminally through a body passageway. Traditionally, stents are made from a metal or other synthetic material with a series of radial openings extending through the support structure of the stent to facilitate compression and expansion of the stent. Although stents may be made from many types of materials, including non-metallic materials, common examples of metallic materials that may be used to make stents include stainless steel, nitinol, cobalt-chrome alloys, amorphous metals, tantalum, platinum, gold and titanium. Typically, stents are implanted within a passageway by positioning the stent within the area to be treated and then expanding the stent from a compressed diameter to an expanded diameter. The ability of the stent to expand from a compressed diameter makes it possible to thread the stent to the area to be treated through various narrow body passageways while the stent is in the compressed diameter. Once the stent has been positioned and expanded at the area to be treated, the tubular support structure of the stent contacts and radially supports the inner wall of the passageway. As a result, the implanted stent mechanically prevents the passageway from narrowing and keeps the passageway open to facilitate fluid flow through the passageway.\nStents can generally be characterized as either balloon-expandable or self-expanding. Traditionally, balloon-expandable stents have been used most often in coronary vessels than in peripheral vessels because of the deformable nature of these stents. One reason for this is that peripheral vessels tend to experience frequent traumas from external sources (e.g., impacts to a person's arms, legs, etc.) which are transmitted through the body's tissues to the vessel. In the case of peripheral vessels, there is an increased risk that an external trauma could cause a balloon-expandable stent to plastically deform in unexpected ways with potentially severe and/or catastrophic results. However, in the case of coronary vessels, this risk is minimal since coronary vessels rarely experience traumas transmitted from external sources.\nSelf-expanding stents are increasingly used and accepted by physicians for treating a variety of ailments. Self-expanding stents are usually made of shape memory materials or other elastic materials that act like a spring. Typical metals used in this type of stent include nitinol and 304 stainless steel. A common procedure for implanting a self-expanding stent involves a two-step process. First, the narrowed vessel portion to be treated is dilated with a balloon but without a stent mounted on the balloon. Second, a stent is implanted into the dilated vessel portion. To facilitate stent implantation, the stent is installed on the end of an inner catheter in a compressed, small diameter state and is usually retained in the small diameter by inserting the stent into an outer sheath at the end of the catheter. The stent is then guided to the balloon-dilated portion and is released from the inner catheter by pulling the outer sheath away from the stent. Once released from the outer sheath, the stent radially springs outward to an expanded diameter until the stent contacts and presses against the vessel wall. Traditionally, self-expanding stents have been more commonly used in peripheral vessels than in coronary vessels due to the shape memory characteristic of the metals that are used in these stents. One advantage of self-expanding stents for peripheral vessels is that traumas from external sources do not permanently deform the stent. Instead, the stent may temporarily deform during an unusually harsh trauma but will spring back to its expanded state once the trauma is relieved. Self-expanding stents, however, are often considered to be less preferred for coronary vessels as compared to balloon-expandable stents. One reason for this is that balloon-expandable stents can be precisely sized to a particular vessel diameter and shape since the ductile metal that is used can be plastically deformed to a desired size and shape. In contrast, self-expanding stents are designed with a particular expansible range. Thus, after being implanted, self-expanding stents continue to exert pressure against the vessel wall.\nCommonly, delivery systems for self-expanding medical devices have a handle arrangement that remains outside of the patient's body during the deployment procedure. One portion of the handle is typically connected to an inner catheter upon which the self-expanding medical device is mounted, and another portion of the handle is typically connected to an outer sheath which restrains the self-expanding medical device in the compressed state. When the distal end of the delivery system is positioned within the patient's body at the intended treatment site, the physician actuates the handle by moving the two portions relative to each other so that the outer sheath is withdrawn from the self-expanding medical device and inner catheter. As a result, the self-expanding medical device expands outward away from the inner catheter. The handle may then be pulled by the physician to withdraw the inner catheter and outer sheath from the patient's body, while leaving the self-expanding medical device implanted in the body.\nPrecise placement of self-expanding medical devices is a concern in most medical procedures. However, precise placement can be more difficult with certain delivery systems due to their design, shape and other factors. Precise placement of self-expanding medical devices is generally a function of the relative movement and placement between the delivery system handle and the patient's body, and the relative movement between the portions of the handle connected to the inner catheter and outer sheath during deployment. A lack of control over any part of this system can result in inaccurate placement of a self-expanding medical device, and thus, less than desirable treatment of the medical condition being treated.\nAccordingly, the inventor believes it would be desirable to provide a new delivery system for self-expanding medical devices."}
{"text": "The present invention relates to a new and distinct variety of rose plant of the hybrid tea class which was originated by me by crossing the variety Baccara, U.S. Plant Pat. No. 1,367 with pollen of the variety South Seas, U.S. Plant Pat. No. 2,184.\nThe primary objective of this breeding was to produce a new rose variety of the hybrid tea class with the color similar to that of Baccara, primarily a greenhouse cut flower variety, but with larger blooms and better formed blooms. This objective was substantially achieved along with other desirable improvements as evidenced by the following unique combination of characteristics which are outstanding in this new variety and which distinguish it from its parents as well as from all other varieties of which I am aware:\n1. A very vigorous, upright, well branched plant.\n2. Most flowers on the flowering stems are borne singly with few side buds.\n3. Fewer petals than its parent, Baccara.\n4. Flower color in overall appearance is near Signal Red with a pronounced yellow area at point of attachment.\n5. Very little fragrance.\n6. Nearly round leaflets, similar in shape to those of Baccara.\nAsexual reproduction of this new cultivar by budding as performed at Wasco, Calif., shows that the foregoing characteristics and distinctions come true to form and are established and transmitted through succeeding propagations."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of locking rings and more particularly to the field of locking rings for electrical energy measuring devices.\n2. Background Art\nElectrical energy measuring devices, both single and poly-phase, are used to indicate the actual electrical power consumption by a particular customer and are well known. One of the most common of such measuring devices is the electro-mechanical induction-type watt-hour meter which is used to measure actual electrical power consumed by most residential customers. The meter is electrically connected in the electrical power supply wiring to the customer's electrical load wiring network and is located at some convenient, accessible location on the customer's premises for minimum disruption of the customers activity. Periodically the meter is read and the customer is billed or charged for the actual consumption of electrical power.\nBecause the meter is located on the customer premises for convenient access by a meter reader, it is also accessible by others and thereby subject to possible theft or tampering. Because, for ease and speed of installation, the meter is normally operably mounted by plugging into a mounting base, it is also easy to remove. While theft of a meter by a determined individual probably cannot be economically avoided, some tampering and vandalism may be avoided by a suitable securing locking ring. Such locking ring also serves to operably connect the meter with the mounting base for maintaining proper electrical contact between the electrical bayonet connections of the meter and the mounting base.\nIn the past such meter locking rings have been formed of a gapped substantially circular metal retainer having a concave securing recess which extended about substantially the entire periphery of the engaged circular flanges of the mounting base and the meter. The metal locking retainer was sufficiently elastically deformable to enable concentric installation and then maintain the adjacent flanges in continuous peripheral sealing contact to exclude dirt and moisture from the mounting base. A soft metal end securing seal clamp is then used to close and seal the gapped ends of the installed retainer ring. Should the end clamp be removed, the impressed soft metal seal is broken and tampering is immediately apparent to the person reading the meter. In the event of an indication of tampering, a false meter reading is immediately suspected and a potential unauthorized power consumption problem identified for proper resolution.\nThe prior art metal securing clamps are expensive, are subject to corrosion and requires the use of a special tool to emboss or form the tattle-tale tampering seal in the soft metal of the end seal clamp."}
{"text": "Parts formed with several portions are known, for example profiled elements of a turbine engine comprising a main body in a composite material and a metal reinforcement. Upon using the profiled element in operation or during production phases, the reinforcements may be worn or damaged while the main body is protected. Thus, it may be advantageous, in maintenance and in production, to replace the reinforcement while retaining the original main body.\nIn order not to weaken the main body, it is preferable to reuse the attachment elements (for example a riveting, tapping hole, etc.) already provided on the main body. Such an operation however requires exact localization, on the new reinforcement, of the positions of these attachment elements. Therefore there exists a need for a device or a method for localizing a position on a part."}
{"text": "In media distribution over an IP network, such as video and TV distribution over the Internet, the bandwidth to a client device will vary depending on various circumstances. When accessing the distributed media content over a mobile data network or a Wi-Fi network the capacity is shared between client devices. Further, individual client devices might enter locations with weaker or stronger signal affecting the bandwidth received by the client.\nToday, variation in bandwidth as conceived by the client device is typically handled by three mechanisms: congestion control mechanisms of TCP (transmission control protocol of the TCP/IP protocol stack), buffering, and adjusting the video bitrate (ABR). Basically, congestion control is handled by the TCP protocol stack which adjusts the retransmission rate of lost packets to adapt the client device to use a fair share of the available bandwidth in the network (or actually in the bottleneck of the transmission). In such system, the client device needs to buffer data since it is not certain the network can offer enough bandwidth required by the video stream. In order to maintain the viewer experience the client device needs to have video data to present, so buffering is needed to absorb variances in bitrate introduced by the network, specifically by the TCP congestion avoidance mechanisms and the jitter introduced by the network. As the capacity and jitter varies in the IP network, the receiving client device must pause the presentation of the current video on the screen to accumulate more video data in its buffer. The accumulation is one method, meaning that the delay will increase and not decrease for a specific session. Such adjustments will introduce delay and the video cannot be considered live distribution due to the added delay.\nAccording to the European Broadcasting Union, EBU, which defines TV standards in Europe, live TV is defined as a broadcasting delay from the ingress to the client device display or screen lower than seven seconds. Delay caused by the adjustments above may however end up in several minutes of delay.\nMore specifically, one of the most common ways to distribute video over the Internet is to use HLS (HTTP Live Streaming) or MPEG-DASH where the video stream is divided into typically 10 second (2-10 sec) video files (segments) making the linear video stream a series of 10 second video files. Every stream is typically represented by several bitrates (different video qualities), each being segmented into equivalent segment files. The client device then requests these files using normal http technology. To ensure that the client device always has video data to present, at least 3 time wise consecutive files are buffered in the device. This means that buffering will impose at least 30 seconds of delay. At start-up the buffer is filled up to a certain level, typically 30 seconds, corresponding to three 10 second segments. If packets are lost, the transport protocol TCP used by the HTTP protocol requests the data again and if uncertain if it can recover the whole segment file, also reduces the bitrate on outgoing traffic by requesting the next corresponding segment file of a lower bitrate to avoid congestion.\nOver IP networks the bandwidth of a distributed video, i.e. the encoded bitrate of the video, is adjusted to the bandwidth available in the network to the client device. The video is encoded in different predefined bitrates (i.e. the level of compression of the video is differentiated, which in turn provides different quality levels of the video). This is typically done by a transcoding system which takes in an encoded video stream and then “re-encodes” it into one or several video streams with different bitrates, qualities and formats for different devices. The network system decides which encoded bitrate is applicable for the specific moment and selects the most suitable video quality to transfer to the client. In existing solutions, ABR adjustments, i.e. adjustment of the video bitrate/quality, are done based on the fill level of the buffer of the client device. More particularly, the decision to change to another encoded bitrate is done by monitoring the buffer fill level in the client device. A decreasing fill level indicates that the bandwidth of the network is lower than needed and the system needs to select a video with lower quality. Shifting to a higher quality video stream is done by simply testing a higher video bit rate and watching if the buffer fill level decreases. If it decreases, the system needs to go back to a lower video bit rate. This means that there will be continuous changes in ABR levels during the operation with the exception if the system is running on the highest ABR level and the buffer is not becoming empty. To avoid glitches in the video when changing between bit streams of different bit rates, the change is performed between segments, i.e. one start using the corresponding new segment of the new bitrate stream.\nAlthough this approach to perform ABR adjustments of video by monitoring the buffer fill level may be applicable for video on demand services, there is a need for an improved method to perform ABR adjustments for video in IP network systems since the buffering of multiple segments at the device and possibly also in edge caches, and also the continuous adjustments of the ABR level increases the buffer fill level at the device, thereby increasing the delay before the video is displayed in the device. The accumulated delay happens because imperfections of using TCP to determine the network bit rate (which continuously is changing) force the client device to buffer more video data to avoid the buffer to run empty which in turn disrupts the video presentation. Typically, today's OTT systems start with an initial delay of 30-60 seconds depending on encoding delay, segment sizes, distribution network and client player implementation. This accumulates over time and e.g., the HLS protocol allows a client to buffer up to 15 minutes."}
{"text": "This invention relates to a pattern synchronous circuit for sorting a parallel signal so that a signal evaluation can be made to a signal after making a serial-parallel conversion of an inputted evaluation pattern in a pattern evaluation apparatus for evaluating a frame signal.\nA frame signal having a frame identification pattern is used in a transmission device test.\nIn this kind of test, a pattern based on actual transmission standards is outputted from a pulse pattern generator and is inputted to an evaluation device unit and the output is observed by a pattern evaluation apparatus or an oscilloscope.\nThe pattern evaluation apparatus normally makes a serial-parallel conversion to a received signal and reduces the signal to a signal speed easy to process and conducts the evaluation. However, timing of inputting the received signal is not determined, so that a leading position of a frame is not determined in a signal after the serial-parallel conversion. As a result of that, it must be constructed so that subsequent evaluations can be made by detecting the frame identification pattern and sorting a parallel signal.\nFIG. 11 shows a configuration of a pattern n apparatus after the serial-parallel conversion. In FIG. 11, numeral 1 is a parallel signal input terminal to which a signal changing a received signal into n parallel signals is inputted, and numeral 2 is a clock input terminal, and numeral 31 is a branch circuit, and numeral 32 is a frame detection circuit, and numeral 33 is a shift circuit, and numeral 34 is a measurement circuit, and numeral 35 is a frame position signal, and numeral 300 is a pattern synchronous circuit.\nThe signal changing the received signal into the n parallel signals is inputted to the frame detection circuit 32 and the shift circuit 33 through the branch circuit 31.\nA clock signal for processing a parallel signal is inputted to the clock input terminal 2. In the frame detection circuit 32, a parallel input signal is monitored and a frame identification pattern is detected and a signal indicating its place is outputted to the frame position signal 35. The shift circuit 33 is a sort circuit for sorting the parallel input signal and in the shift circuit 33, the parallel input signal is sorted by the frame position signal 35 and in the measurement circuit 34, an output is produced so that it becomes a measurable state, namely the lead of a frame becomes a leading position of the parallel signal.\nIn this case, the shift circuit 33 can implement a pattern synchronous operation by performing processing as described in a truth table shown in FIG. 12, for example, when the number n of parallel signals is 16. The truth table shown in FIG. 12 indicates that SEL corresponds to a frame position signal and a signal is selected from the parallel signals of 16 bits as shown in output by SEL input and is outputted. Also, the places described as A(+1), B(+1), C(+1), . . . indicate signals after one clock of inputs A, B, C, . . . ,\nHere, when the shift circuit 33 of FIG. 11 produces an increase in circuit scale, particularly an increase in the number n of signals, integration into one integrated circuit becomes difficult. Thus, the shift circuit 33 is implemented by a plurality of integrated circuits.\nNext, a pattern synchronous circuit according to a conventional art will be described using FIG. 7. In FIG. 7, numeral 1 is a parallel signal input terminal, and numeral 2 is a clock input terminal, and numeral 3 is a parallel signal output terminal, and numeral 21 is a branch circuit, and numeral 22 is a frame detection circuit, and numeral 23 is a shift part, and numerals 24a to 24d are shift circuits A to D constructing the shift part 23, and numeral 25 is a frame position signal. The shift part 23 comprises a plurality of the shift circuits 24a to 24d. As a result of that, the shift circuit per circuit performs sorting of the number n of parallel signals and data processing of the number divided by the number m of integrated circuits.\nA signal inputted to the parallel signal input terminal 1 is branched to the frame detection circuit 22 and all the shift circuits 24a to 24d through the branch circuit 21. The frame detection circuit 22 detects a frame identification pattern in the parallel signals and outputs a signal to the frame position signal 25. The frame position signal 25 is connected to all the shift circuits 24a to 24d, and the shift circuits 24a to 24d sort the parallel signals by the frame position signal 25.\nAs examples, FIGS. 8 and 9 indicate truth tables showing operations of the shift circuit A and the shift circuit B where the number n of parallel signals is 16 bits and four shift circuits are used.\nWhen the frame position signal is 0, the shift circuit A outputs signals A, B, C, D inputted to D0 to D3. At that time, the shift circuit B outputs signals E, F, G, H inputted to D4 to D7 since the shift circuit B outputs signals subsequent to the signals outputted by the shift circuit 24a. Also, when the frame position signal is 1, the shift circuit 24a outputs signals B, C, D, E inputted to D1 to D4, and the shift circuit 24b outputs signals F, G, H, I which are signals subsequent to the signals outputted by the shift circuit 24a. Truth tables of operations of the shift circuit 24c and the shift circuit 24d are omitted, but the truth tables can be derived in a manner similar to truth values of FIGS. 8 and 9.\nOperations of a pattern synchronous circuit 200 of FIG. 7 will be described concretely using FIGS. 10A and 10B.\nIt is assumed that the number of parallel signal inputs is 16 bits and for the purpose of illustration, parallel signals of 16 bits are described as a, b, c, . . . , o, p and the leading of a frame begins with a. Expressions such as p(−1) indicate a bit earlier than p by one clock.\nFIGS. 10A and 10B show an example in which a frame of an input signal begins with BIT 7 of the parallel signal input terminal. The frame detection circuit 22 outputs a value 7 to the frame position signal 25 when detecting the frame in beginning with BIT 7. Then, the shift circuit 24a selects BITs 7 to 10 from the parallel signals of 16 bits, and also the shift circuit 24b selects BITs 11 to 13 and outputs them. Similarly, the shift circuits 24c and 24d select respectively corresponding 4 bits from input signals and thereby, an output signal in which the leading of the frame begins with BIT 0 of the parallel signals as a whole can be obtained.\nHere, all the information of the inputted parallel signals needs to be inputted to all the shift circuits 24a to 24d in order to cope with a start position of the frame in all the states. For that purpose, all the bits of the input signals are branched to all the shift circuits 24a to 24d by the branch circuit. Then, all the shift circuits 24a to 24d require the number of input I/Os corresponding to the number n of parallel bits of the input signals.\nAs described above, in the pattern synchronous circuit according to the conventional art, all the bits of the parallel signals inputted to the parallel signal input terminal 1 are branched by the branch circuit 21 and are inputted to all the shift circuits 24a to 24d constructing the shift part 23. Then, the respective shift circuits 24a to 24d select only the number of signals obtained by dividing the number n of parallel signals of the output signals by the number of sort integrated circuits from the input signals and output it on the basis of the frame position signal 25 outputted by the frame detection circuit 22. An output signal in which the frame begins with BIT 0 of the parallel signals can be obtained by merging outputs of all the shift circuits 24a to 24d. \nHowever, in the conventional pattern synchronous circuit described above, an input signal is branched to all the shift circuits, so that the number of wirings increases in proportion to the number n of parallel signals and the number m of shift circuits. Also, since all the bits of a parallel signal input are inputted to the shift circuits, the number of I/Os of the shift circuits becomes large with an increase in the number n of parallel signals and thus a large package is required. As a result of that, miniaturization of the pattern synchronous circuit is difficult."}
{"text": "Various types of such roll-off pistons are known in the prior art. Firstly, relatively heavy roll-off pistons exist which are made of steel sheet metal and which have a fully utilizable interior volume and an end stop, that is, a so-called buffer support, specifically a support for a stop buffer that acts substantially in the axial direction of the air spring. The roll-off piston made of steel sheet metal is produced as a deep-drawn part with a conical sealing seat for receiving the bellows, and has a welded-in base part with welded-on support tube. The entire component is accordingly heavy and is expensive to manufacture.\nSecondly, there exist relatively lightweight, single-piece plastics pistons which have an unutilized or only partially utilized interior volume and an end stop (buffer support). A disadvantage of such constructions is a loss of comfort resulting from a small utilizable piston air volume. The suspension characteristic is then relatively hard.\nFurthermore, the prior art also discloses two-part plastics pistons with a utilizable interior volume without buffer support.\nEP 1 862 335 B1 discloses for this purpose an air spring piston which is made of a cup-shaped part and of a cover part, the parts being butt-welded in the region of the walls thereof. Here, the cup-shaped part has a base wall which has an insert for connecting to the vehicle axle. A disadvantage here is the lack of an end stop (buffer support).\nUnited States patent application publication 2010/0127438 A1 discloses a plunger, in the form of a hollow body, for an air spring, the plunger being composed of two parts connected to one another in an air-tight manner, specifically of a cup-shaped lower part with a base and a shell and of an upper part. In this publication, too, no end stop is provided."}
{"text": "1. Field of the Invention\nThe present invention relates to an imaging device.\nPriority is claimed on Japanese Patent Application No. 2012-208597, filed Sep. 21, 2012, the contents of which are incorporated herein by reference.\n2. Description of Related Art\nIn recent years, as a method of diagnosing cancer, a method using a fluorescent pigment which is typified by indocyanine green (ICG) has been actively studied. A technique is known which forms a filter using Fabry-Perot interference on an image sensor using a general CMOS process to manufacture an optical element for detecting light with a specific wavelength, such as fluorescence (for example, see Non-patent Document: “Nicolaas Tack, three others, A Compact, High-speed and Low-cost Hyperspectral Imager, SPIE, 2012”). Hereinafter, a pixel in which a Fabry-Perot filter is formed on a photodiode is referred to as a Fabry-Perot filter pixel (FPF pixel).\nIn addition, a technique is known which applies the Fabry-Perot filter to an endoscope to capture a fluorescent image, thereby separately capturing the image of a cancer tissue part using a fluorescent pigment (for example, see PCT International Publication No. WO2009/072177). Furthermore, a technique is known which generates an RGB image that is generally captured by an imaging device, such as a digital camera, from a fluorescent image to capture a visible image sensed by the human eye (for example, see Japanese Patent No. 3781927).\nThe technique disclosed in Non-patent Document: “Nicolaas Tack, three others, A Compact, High-speed and Low-cost Hyperspectral Imager, SPIE, 2012” will be described with reference to FIGS. 6 and 7. FIG. 6 is a schematic diagram illustrating an example of an FPF pixel according to the related art in which a Fabry-Perot filter is formed on a pixel of a CMOS image sensor. In the example shown in FIG. 6, a plurality of pixels F1 to F8 are formed on a semiconductor substrate 5-5. A photodiode (PD) 5-4 is formed in each of the pixels F1 to F8. In addition, dielectric layers 5-1 and 5-3 and an interlayer film 5-2 are formed, as the Fabry-Perot filter, on a light-receiving surface of each of the pixels F1 to F8. The interlayer film 5-2 which is formed above each of the pixels F1 to F8 is interposed between the dielectric layers 5-1 and 5-3. The thicknesses of the interlayer films 5-2 formed above the pixels F1 to F8 are different from each other such that light components in different wavelength bands are incident on the photodiodes 5-4 provided in the pixels F1 to F8. According to this structure, the pixels F1 to F8 can detect, for example, light in a narrow transmission band, as shown in FIG. 7.\nFIG. 7 is a graph illustrating the wavelength of light passing through the Fabry-Perot filter according to the related art. In the graph shown in FIG. 7, the horizontal axis is a wavelength and the vertical axis is transmittance. In the example shown in FIG. 7, the Fabry-Perot filter (the dielectric layers 5-1 and 5-3 and the interlayer film 5-2) formed on the pixel F1 transmits light with a narrow band wavelength of about 420. Therefore, the pixel F1 can detect light with a narrow band wavelength of about 420 nm. Similarly to the pixel F1, the pixels F2 to F8 can detect light with a narrow band wavelength, depending on light with wavelengths transmitted through the farmed dielectric layers 5-1 and 5-3 and the farmed interlayer film 5-2. The band of light transmitted through the dielectric layers 5-1 and 5-3 and the interlayer film 5-2 formed on each of the pixels F2 to F8 is as shown in FIG. 7.\nA fluorescent image can be captured by the pixels F1 to F8 shown in FIG. 6 and an RGB image can be generated based on the fluorescent image, as described in the related art. For example, in the technique disclosed in Japanese Patent No. 3781927, the spectral characteristics of blue light are generated based on the values detected by the pixels F1 to F3, the spectral characteristics of green light are generated based on the values detected by the pixels F4 and F5, and the spectral characteristics of red light are generated based on the values detected by the pixels F6 to F8. In this way, it is possible to acquire an RGB image.\nFor example, in the actual surgical operation, it is necessary to capture images of a part separated by a fluorescent pigment and the other parts at the same time and present the images to a doctor. Therefore, it is necessary to attach a device capable of capturing both an RGB image and a fluorescent image to the leading end of the endoscope. In addition, it is necessary to reduce the size of the leading end of the endoscope. The use of the above-mentioned technique makes it possible to generate both a fluorescent image with narrow-band sensitivity characteristics and an RGB image while reducing the size of the leading end of the endoscope."}
{"text": "1. Field\nThe invention relates to the technical field of vehicle tires and, in particular, to that of the manufacture of tires comprising fibres injected in a given direction into one of the rubber components of which said tire is made.\n2. Description of Related Art\nPublication EP 1 213 130 provides a specific idea of this type of tire, in which short reinforcing fibres, or inserts, are embedded in the tread in an oriented fashion so as to improve the mechanical properties of this part of the tire in a given direction.\nPublication EP 1 213 383 proposes a means for embedding said inserts into a rubber component intended to be assembled to form a tire tread structure.\nThis device comprises a hollow needle capable of penetrating the rubber component, and inside which the reinforcing fibre runs, a moving support capable of performing a reciprocating movement and supporting said needle, a cutting blade for cutting off the fibre to be embedded, and at least two clamping means for maintaining and adjusting the position of the fibre. However, these clamping means have the effect of placing the reinforcing fibre under compression.\nApplication PCT/EP2010/062907, not published, describes an apparatus and a method for overcoming this compression problem.\nThis device for embedding an insert into a rubber profiled element in a given direction comprises:                a means for feeding the thread continuously,        a moving support capable of moving on a structure in a rectilinear reciprocating movement between a raised position and a lowered position,        a hollow needle comprising a channel inside which said thread is able to run, which needle is fixed by one of its ends to the base of the moving support and comprising at its other end a tip capable of penetrating said rubber profiled element,        a cutting means arranged on the structure, comprising a blade capable of cutting off the thread on each cycle of the reciprocating movement of the needle, when said tip is in a raised position so as to release a given length of thread constituting a tail at the outlet of the needle.        \nStill according to this publication, there is described a method of operating the device which comprises the following steps:                A—when the needle is in a raised position above the surface of the profiled element, a given length of thread is released on the tip side of said needle so as to constitute a tail,        B—the needle is lowered so as to cause the tip of said needle to penetrate the profiled element to a given depth,        C—the needle is extracted from the profiled element while allowing the thread to run inside the needle,        D—the tip of the needle is raised by a given height above the surface of the profiled element so as to release a length of thread between the tip of the needle and the surface of the profiled element,        E—the thread is cut off at the surface of the profiled element so as to release a new tail of given length.        \nOn each cycle, the tail situated outside the needle on the tip side is driven into the profiled element. As the needle is raised again, the tail is held in position in the bottom of the perforation produced in the rubbery material which grips it, and which prevents the thread contained in the needle from rising back up, such that the raising movement allows the thread to be made to run along inside the needle.\nDuring this movement, the portion of thread contained between the inlet of the needle and the tail is kept under tension. The thread is pulled through the needle rather than pushed as was proposed in the apparatus forming the subject matter of publication EP 1 213 383.\nHowever, it is observed that, to make the needle penetrate the rubber profiled element, it is necessary to develop considerable forces which are not without consequence on the dimensioning of the needle and for the mechanisms to be employed to bring about its reciprocating movement. The effect of these forces is to deform the surface of the profiled element at the point of penetration of the needle, and they can moreover cause the needle to break, in particular when the angle of penetration of the needle into the profiled element is less than 90°."}
{"text": "Remote Procedure Call (RPC) is a mechanism by which an application residing on one machine requests the services of an application residing on another machine. A machine (as the term is used here and in the appended claims) may be a physical CPU or computer, a virtual CPU or virtual computer (otherwise known as a virtual machine or VM), or it may be a space of addressable memory in a system that isolates applications or processes so that they may not directly invoke one another. RPC requires one application to send one or more messages to another application in order to invoke a procedure of the other application. The recipient application may reply by sending one or more messages back to the requesting application. The term “RPC” identifies both the mechanism by which this occurs and the instance of such an exchange between applications.\nApplications engaged in RPC conventionally do not message each other directly. Instead they rely on intervening software called “middleware” to translate between the wire protocol and the programming-language-specific functions or class methods of the application. This application uses the term “service” to generically refer to a function (including a procedure or a class method) or any action to be taken or performed, including access to web sites and databases. The term “wire protocol” refers to the transfer-protocol/message-encoding combination that the programs use as their common language for communication. The transfer-protocol is the mechanism by which the message transfers between the applications, and the message-encoding is the format in which the message is represented. The Object Management Group's IIOP (OMG Internet Inter-ORB Protocol) and the Open Software Foundation's DCE (OSF Distributed Computing Environment) are examples of wire protocols, and each defines both a transfer-protocol and a message-encoding. RPC has in the past been implemented using IIOP, DCE, or proprietary protocols. HTTP is an example of a transfer-protocol that is independent of message-encoding. The IIOP and DCE protocols have themselves been encoded into messages that tunnel through HTTP. ( )\nEach message-encoding is a particular format for representing a message. Any set of data may be represented in a particular format, not just RPC messages. This application uses the term ‘data-package’ to refer to each such set of data, and it uses the term ‘encoding’ to refer to the format in which a data-package is represented. An RPC message is a kind of data-package, and a message-encoding is the encoding of an RPC message. A tab-delimited file is an example of a simple encoding. The binary representation of a C programming language data structure is an example of a more complex encoding. A data-package may be thought of as a collection of individual data items. In a tab-delimited file, each row is a data item, as is each field in each row. In a C data structure, each field is a data item, as is each nested data structure and each field in each nested data structure. Notice that a data item may contain one or more other data items.\nAn ‘argument’ is a particular kind of data item. The term argument is used in the software industry to refer to a data item that is passed to a function or that is returned from a function. When the signature of a function is represented in the Java programming language, the arguments of the function consist of the comma-delimited terms that appear in parentheses and the term that represents the return value of the function. Consider the following function signature:\nEmployeeID addEmployee(Name n, Address a, JobType t)This function has four arguments whose classes are EmployeeID, Name, Address, and JobType. The Address class might be defined as follows:\nclass Address {  String street;  String city;  String state;  String zip;}In other contexts one might interpret the street, city, state, and zip data items to each be an argument of the function, but this application defines ‘argument’ such that these data items are not arguments. The arguments of a function are those data items that the signature of the function identifies, where for purposes of this definition the definitions of the argument types are not considered to be part of the signature.\nA function is only one kind of service, and XML RPC messages can exist in the absence of functions. Hence, it is necessary to define the term ‘argument’ in terms of the message encoding. Therefore, given a data item X contained in message M1 expressed in encoding E, X is an argument when it satisfies both of the following requirements: (1) there exists a message M2 also expressed in E such that M2 is constructed by adding a data item Y to M1 where neither X contains Y nor Y contains X; and (2) X is not contained in any data item having property (1). It is possible to define an encoding that hard-codes arguments relative to another encoding so that the defined encoding expresses a subset of the messages of the other encoding. Consequently, a message cannot be examined for arguments without interpreting the message with respect to a particular encoding.\nEach argument is an input argument, an output argument, or both an input and an output argument. The return value is always an output argument, but function signatures usually do not indicate which if the remaining arguments are input arguments or output arguments. The input arguments and output arguments must be established by convention. The input arguments are those arguments that the software that invokes the function uses to pass data items to the function. The output arguments are those arguments that the function uses to pass data items to the software that invoked the function. It is possible for one argument to be used for both purposes. In this case, the argument appears in both the RPC request message and the RPC reply message, although the values of the occurrences may differ. The OMG CORBA specification provides further treatment of the distinction between input and output arguments in the context of RPC.\nIn addition to containing data items, a data-package may contain labels. A label is an identifier that provides information about one or more data items. A label that provides information about a data item is said to be a label for the data item. No data item is a label, since the encoding assumes responsibility for providing the association between the data item and the label.\nThis application defines two kinds of labels: type labels and semantic labels. A type label identifies the data type of each associated data item. Data types include the primitive data types that programming languages define. The primitive data types include integers, floats, longs, strings, and booleans. According to this definition, data types also include primitive data structures, such as arrays, records, and vectors (e.g. the Java Vector class). Data types provide the information that is necessary to represent the data item in memory so that a programming language may use the data item. However, there is no requirement that any two programming languages or any two implementations of RPC represent the same data types in the same way.\nA semantic label is an identifier that names an associated data item according to the data item's role within the message. The purpose is to allow software that consumes the message to locate data items by role. The semantic label of an argument is analogous to the name that a method signature gives to the argument. Again consider the following signature:\nEmployeeID addEmployee(Name n, Address a, JobType t)This function has arguments named ‘n’, ‘a’, and ‘t’. It also has an unnamed argument whose type is EmployeeID. The names may be used as semantic labels for their associated arguments. The semantic label of a field found in a class, record, or structure is analogous to the name given to that field. Consider the previously given definition of the Address class. It has four data items. Each has a data type of string, and their names are ‘street’, ‘city’, ‘state’, and ‘zip’. These names may be used as the semantic labels for the associated data items. However, unlike the names of a programming language structure, multiple data items may have the same semantic labels. For example, two data items may each have the ‘street’ label to indicate that they each represent one line of the street address. On the other hand, an array or a vector may contain multiple data items such that a semantic label is associated with the entire array or vector but not with each data item.\nThis application classifies encodings as either self-describing or non-self-describing. A self-describing encoding associates a semantic label with each argument and includes both the label and the association within each data-package. A completely self-describing encoding associates a semantic label with every data item except possibly for the data item constituents of arrays and vectors. A semantic label allows a program to identify a data item without requiring the program to know in advance the location of the data item within the data-package and without requiring the program to know the exact structure of the data-package. One example of a data-package that uses a self-describing encoding is a tab-delimited file that begins with a row consisting of the names of each field in a row. Another self-describing encoding is the binary representation of a hash table, which assigns a name (or key) to each data item in the table. Non-self-describing encodings do not contain semantic labels for their arguments. A program that reads a data-package expressed in a non-self-describing encoding must assign semantics to each argument by applying information found in the application and not in the data-package. For example, a program that reads a tab-delimited file to load records into a database must know in advance how to correlate the fields in the file to the columns of the database table.\nThis application also classifies encodings as either specific or generic. A specific encoding is an encoding that is designed to represent an application-specific or industry-specific data structure. A generic encoding is an encoding that is designed to represent any or nearly any data structure, regardless of application or industry. For example, the ASCII form of a database report is a specific encoding since it can express the data found in a particular database but not the data found in any database. Packet headers in communications protocols also use specific encodings. An implementation of the C programming language uses a generic encoding to represent in-memory data structures; this is a generic encoding because one encoding scheme is capable of representing any C data structure.\nFinally, this application distinguishes between binary encodings and text-based encodings. A text-based encoding is a human-readable and human-writable encoding. Data-packages expressed in a text-based encoding can be loaded into a text editor and can be read and modified by humans using the text editor. A human can also create a data-package expressed in a text-based encoding by writing the message completely from scratch using the text editor. Text editors include vi, emacs, MS-DOS EDIT.EXE and Microsoft Notepad. Binary encodings are those encodings that do not satisfy these human-accessibility criteria. Tab-delimited files are examples of data expressed in a text-based encoding, while the in-memory representation of C data structures is an example of a binary encoding.\nThe IIOP and DCE wire protocols use non-self-describing generic binary encodings. Of all the example encodings discussed in this application, the IIOP and DCE encodings are most closely analogous to the in-memory binary representations of C or C++ data structures. However, the analogy is not perfect because the encodings have been designed to support the data structures found in a variety of languages, and because they have been designed to support the translation of data between a variety of platforms.\nAlthough IIOP and DCE are non-self-describing, it is possible to define services in IIOP and DCE that accept self-describing parameters and that invoke other services using the identified parameters. Microsoft's COM (Component Object Model) accomplishes this on top of DCE through the set of services found in its IDispatch interface. By invoking the services of IDispatch, an application can emulate the functionality that would be available through a wire protocol that is natively self-describing. However, DCE remains non-self-describing, since the services of IDispatch are expressed in the same non-self-describing encoding in which all other service invocations of the protocol are expressed. IDispatch effectively defines a high level protocol that is layered on top of DCE. More work is required to develop software that invokes self-describing services through IDispatch than is required to develop software that invokes the native non-self-describing services of DCE. Additional effort is required because to invoke a single self-describing service the developer must write software that conforms to the IDispatch protocol. Invoking a single non-self-describing service only requires the developer to invoke a single function in the code that DCE IDL compilers automatically generate for the developer.\nIIOP defines a data-type called ‘ANY’. A data item expressed in this type is a container for other data items. The ANY data type associates a type label with each data item that it contains, but it does not associate a semantic label with any of these data items. It also does not associate a type label with the instance of the ANY data item itself, unless the instance happens to be contained within another ANY data item. Consequently, an implementation of IIOP that receives a message containing an instance of the ANY data type cannot determine that the instance is of type ANY simply by examining the message. The recipient of the IIOP message must know in advance that the argument is of type ANY. In fact, the recipient must know the types of all the arguments in order to extract the arguments from the message. This aspect of IIOP makes programs that use IIOP strongly sensitive to changes in the messages. Adding, removing, or changing argument types in the programs that generate the IIOP messages necessitates making analogous changes in all the programs that receive the messages. This would not be necessary in programs that relied on data type labels that the message itself provides.\nIIOP and DCE are also stateful, bi-directional, and complex. These protocols are stateful because the client and server must maintain state information to successfully transfer messages between them. The protocols are bidirectional because the server must sometimes asynchronously send messages to the client. Finally, IIOP and DCE are complex because they are designed to provide the union of the features found in a wide variety of programming languages. This latter property of IIOP and DCE allows any two programs to communicate with service invocations that support the rich variety of invocation mechanisms that are common to the programming languages in which the two programs are written.\nTraditional middleware implementations are based on the IIOP and DCE protocols and consequently suffer from the following drawbacks:                Because the encodings are generic, translation software is required to allow RPC messaging between a program that uses a specific encoding and a program that uses a generic encoding. Translation software is also required between programs that use different specific encodings. If IIOP or DCE is to be the wire protocol, the translation software must be installed between each program that uses a specific encoding and the communications channel that connects the programs. Furthermore, generic encodings are necessarily more removed from the data structures that are common to a given industry, which prevents most industry experts from examining messages, since expertise with the generic encoding would be required in addition to expertise with the industry data structures.        Because the encodings are binary, special software is required to read, modify, or write the RPC messages. If the message is expressed in a non-self-describing encoding, the special software must be application-specific or industry-specific software. These factors limit the accessibility of messages expressed in the IIOP and DCE encodings, since the translation software must exist, and since one must have translation software on hand to access the messages.        Because these encodings are non-self-describing, programs can't extract and selectively process the data items of an arbitrary message without first being specifically configured to recognize the particular kind of message. Programs that do support configured access to non-self-describing messages generally require separate configuration information for each kind of non-self-describing message. As the message-generating programs evolve the messages and add new message types, the configuration information of the other programs must be updated. To further complicate matters, each program may require different configuration details or separate configuration entry procedures. Finally, if a message-generating program changes the format of a non-self-describing message, the message-consuming programs must also change their knowledge of that message, since otherwise the these programs would not be able to extract and identify the data items found in the message.        Because the IIOP and DCE protocols are stateful, the protocols are not well-suited for messaging using HTTP, since HTTP is a stateless protocol. HTTP provides only a limited mechanism for maintaining client state. Under this mechanism, the client maintains the state information, which means that client-provided state information cannot be trusted to belong to the same client over time, unless the client connects through an authenticated connection. One must therefore inventively design a stateful protocol on top of HTTP in order to tunnel IIOP or DCE through HTTP. It is desirable to use HTTP as the transfer protocol because of the ubiquity with which it is installed on operating systems, and because it is typically the only protocol that network firewalls allow to pass.        Because these protocols are bi-directional, the protocols are not well-suited for messaging using HTTP, since HTTP is unidirectional. For example, an IIOP client may register a callback function with an IIOP server. When it comes time for the server to invoke the callback function, the server cannot issue an HTTP request to the client unless the client is also configured to be an HTTP server. Administrative and security requirements normally make it undesirable to configure a client as an HTTP server. Therefore, one must inventively design a bi-directional protocol on top of HTTP in order to support IIOP and DCE bi-directional behavior.        Because the encodings are intended to bridge applications written in almost any two programming languages on almost any two platforms, the encodings are proportionately complex. IIOP and DCE take the superset approach to bridging applications by ensuring that the encoding can represent nearly every kind of service invocation that programming languages are capable of expressing. To create software that uses IIOP or DCE, a software developer must define the signatures of these services. The superset approach complicates the definitions and minimizes the pool of skilled experts that can develop for IIOP and DCE. This approach also complicates the tools that manipulate messages expressed in these encodings.        Because of the complexity of the IIOP and DCE protocols, implementations from two different vendors are generally incompatible. To get two programs to communicate over a network, the computers running the programs must conventionally run middleware software from the same vendor. Vendors have implemented gateway software to bridge between different implementations of these protocols. The CORBA 2.0 specification has mitigated this problem to a large degree, but it has not eliminated it completely. Thus, integrating programs residing in different organizations often requires either both organizations to agree on the installation or one company to buy additional software to provide a gateway between the existing middleware installations.        \nTherefore, there is a need for improved RPC techniques that do not suffer from the above drawbacks. The techniques would provide the following benefits:                It would be beneficial to have an RPC mechanism that allows messages to be expressed in application-specific or industry-specific encodings. This would allow programs to communicate without requiring messages expressed in specific encodings to be translated into messages expressed in generic encodings. Provided that the messages were also text-based, this would also allow industry experts to create and examine messages without requiring that they have expertise in a generic encoding. Text-based messages in specific encodings would also be easier to read than text-based messages in generic encodings, since the generality of generic encodings makes generically-encoded messages more verbose.        It would be beneficial to have an RPC mechanism that uses text-based messages, since special software would not be needed to read, write, or modify the messages. Text editors would suffice. People who have expertise in industry data structures would have less reliance on the expertise of others to translate messages into a form that they can understand.        It would be beneficial to have an RPC mechanism that is based on a self-describing encoding such that all messages are at some level expressed in the self-describing encoding. This approach would simplify the programs that access the data items of the messages, since the programs would no longer require configuration information. Furthermore, because the programs that generate and consume the messages could access data items by label, the programs would be less sensitive to changes in the message types.        It would be beneficial to have an RPC mechanism that operated statelessly over HTTP so that less complex software is required to implement RPC over HTTP.        It would be beneficial to have an RPC mechanism that did not support bi-directional communications so that less complex client software is required to implement RPC over HTTP. A unidirectional RPC would not support asynchronous callbacks, and an application that is based on the RPC would have to either select a synchronous solution or install an HTTP server at both ends of the communication channel.        It would be beneficial to have a generic encoding that precludes some of the less widely used features of IIOP and DCE so that the encoding is simpler than IIOP and DCE. It would be beneficial for this generic encoding to contain only those features that are necessary for most business-to-business electronic commerce over the Internet, thus idealizing the encoding for this domain of applications. The approach maximizes the number of software developers that are qualified to produce solutions with the encoding and minimizes the time required to produce each solution.        It would be beneficial to have an RPC framework whose simplicity allows two programs to engage in RPC communications over a network without requiring software from a single vendor to be installed on both of the computers that are running the programs. Software from different vendors should be compatible by virtue of using the nearly ubiquitous HTTP protocol stack to transmit messages whose encodings use a well-defined and well-accepted text-based syntax.        "}
{"text": "The following includes information that may be useful in understanding the present invention(s). It is not an admission that any of the information provided herein is prior art, or material, to the presently described or claimed inventions, or that any publication or document that is specifically or implicitly referenced is prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to a resist pattern forming method by immersion exposure of a lithography step in a semiconductor manufacturing process, and a manufacturing method of a semiconductor device.\n2. Description of the Related Art\nImmersion exposure is a technology in which a space between the surface of a resist film and a lens of an exposure device is filled with a liquid to perform exposure on the resist film formed on a substrate to be treated. A device for use in such an exposure process is disclosed in, for example, Jpn. Pat. Appln. KOKAI Publication No. 10-303114. In Jpn. Pat. Appln. KOKAI Publication No. 10-303114, there is disclosed a device in which a substrate to be treated is all immerged into water in a stage to which water can be supplied, and the exposure is performed while this stage is relatively moved with respect to the exposure device. In the device having such configuration, the liquid is supplied to the whole stage. Therefore, there is a problem that the liquid overflows from the stage, when the stage is moved at a high speed. The stage cannot be driven at a high speed.\nIn Immersion Lithography; its Potential Performance and Issues, Proc. of SPIE Vol. 5040, pp. 724 to 733 by Soichi Owa and Hiroyuki Nagasaka, as a measure against disturbance of the liquid due to the movement of the stage, there is disclosed a technology to drive the stage while locally supplying the liquid to a portion to be exposed. According to this technology, it is possible to move the stage at a high speed. When such technology to supply the liquid locally is used, water is left in an exposed region and the like of a portion from which the lens is detached. In consequence, when the resist film is exposed and thereafter heated in this state, there is a problem that a watermark is generated, or an abnormality is generated in a resist pattern owing to a temperature drop in the portion where water exists.\nOn the other hand, in New Cover Material Development Status for Immersion Lithography, Web Publication of International Symposium on Immersion and 157 nm Lithography by Keita Ishizuka et al., it is disclosed that components of a resist are eluted from the resist film into the liquid in the immersion exposure. According to this document, such elution can be avoided by disposing a cover film on the surface of the resist film. As to a state of the surface of the cover film, it is preferable that a contact angle of the surface to water is 90° or more. However, in a case where the cover film having such a large contact angle is used, there might occur a problem that the watermark is easily generated on the cover film, and a defect due to the watermark is generated in the resist pattern. In this document, there is also disclosed an example in which a film having a contact angle of 68° or 78° to an immersion liquid (pure water) is used as a hydrophilic cover film."}
{"text": "Audience response systems are employed to retrieve (or receive) responses from a group of individuals at a central location. Such systems may be used in classroom settings, corporate meetings, or in other gatherings of individuals. Wireless audience response systems may include at least one base unit and a plurality of handheld units. Each handheld unit typically includes a keypad for inputting user responses.\nIn one known embodiment, the base unit controls communication by, for example, polling each handheld unit. If a user has entered a response into the keypad, the response is transmitted to the base unit when the handheld unit is polled. After receipt of the response, the base unit may send an acknowledgment signal."}
{"text": "As the value and use of information continues to increase, individuals and businesses seek additional ways to process and store information. One option available to users is information handling systems. An information handling system generally processes, compiles, stores, and/or communicates information or data for business, personal, or other purposes thereby allowing users to take advantage of the value of the information. Because technology and information handling needs and requirements vary between different users or applications, information handling systems may also vary regarding what information is handled, how the information is handled, how much information is processed, stored, or communicated, and how quickly and efficiently the information may be processed, stored, or communicated. The variations in information handling systems allow for information handling systems to be general or configured for a specific user or specific use such as financial transaction processing, airline reservations, enterprise data storage, or global communications. In addition, information handling systems may include a variety of hardware and software components that may be configured to process, store, and communicate information and may include one or more computer systems, data storage systems, and networking systems.\nAn information handling system may include a voltage regulator to provide a constant voltage level and a current to power the system. For example, a voltage regulator may receive an input voltage and produce an output current at a predetermined output voltage required by a load, i.e., the circuit element(s) for which it is providing power. Moreover, modern information handling systems may include components that maintain current requirements across a broad range from relatively high peak currents to very low stable currents. More particularly, voltage regulators may be required to maintain a high efficiency, or low power loss, over such ranges. In particular, a direct current to direct current (DC-DC) voltage regulator may include a controller, one or more drivers, and one or more power stages. Furthermore, a power stage may include one or more metal oxide-semiconductor-field-effect-transistors (MOSFETs), which may be driven by the drivers.\nAdditionally, some voltage regulators may be capable of operating in multiple phases. To this end, the concept of a phase for a voltage regulator may typically refer to combining an inductor and a power stage to form one phase. A power stage (which may also be known as a Pstage) may comprise a high-side MOSFET and a low-side MOSFET each driven by a driver. In some embodiments, a power stage may be fully integrated into a single integrated circuit, known as an integrated power stage.\nA multi-phase voltage regulator may include multiple instances of such phases. Alternatively, a multi-phase voltage regulator may be thought of as a combination of single phase voltage regulators. For example, a multi-phase voltage regulator may include a plurality of single phase voltage regulators coupled in parallel to provide varying ranges of output current. During periods of high loads, the multi-phase voltage regulator may function with all phases in operation. In contrast, for lower loads, it may employ phase-shedding and operate with a reduced number of phases.\nTo manage power delivery and consumption by voltage regulators, power control systems in information handling systems often execute power management algorithms. For such power management algorithms to effectively manage power, accurate power, current, voltage, and/or other measurements must be obtained. An inherent conflict in obtaining power measurements is that measurement circuitry itself may add power losses to a system."}
{"text": "1. Field of the Invention\nThe present invention relates to an electronic component feeding apparatus for conveying electronic components by intermittent movement of a belt.\n2. Description of the Prior Art\nJapanese Patent Application Laid-Open No. Heisei6-232596 has disclosed an apparatus for conveying chip components using a belt.\nThis apparatus has an endless belt wound around a pair of pulleys, a ratchet mechanism for intermittently moving a belt through a predetermined distance by driving one of the pulleys, and a lever mechanism for transmitting the power to the ratchet mechanism. When the power is transmitted to the ratchet mechanism by moving an operating portion of the lever mechanism in a predetermined direction, one pulley is turned through a predetermined angle in a predetermined direction by the ratchet mechanism, and the belt is moved through a predetermined distance in a predetermined direction by the turning of the pulley. Chip components, arranged in an orderly way on the belt, are conveyed in the same direction by the movement of the belt, and the foremost chip component is fed to a predetermined takeout position. The foremost chip component fed to the takeout position is taken out by using a suction nozzle. When the operating portion of the lever mechanism is moved again, the succeeding chip component is fed to the takeout position in the same way as described above.\nThe apparatus of this type has widely been used for a component mounter, component mounting line, and the like. In recent years, there has been demanded a mechanism which can respond to high-speed component takeout cycle, specifically a mechanism which can take out components in a very short cycle as short as 0.1 sec.\nHowever, on the aforesaid apparatus, since the operating portion of the lever mechanism is directly connected to an driving portion of the ratchet mechanism, the speed at which the belt moves is substantially proportional to the speed at which the operating portion of the lever mechanism is moved. That is to say, if the operation speed of the operating portion of the lever mechanism is increased to respond to the high-speed component takeout cycle, the belt movement speed increases substantially proportionally, so that the chip components on the belt cannot follow the belt initial speed. Thereby, a slip occurs between the belt and the chip components on the belt, so that only the belt moves while the chip component is left, with the result that components cannot be conveyed satisfactorily by means of the belt."}
{"text": "Conventional fluted twist drills have the problem of excessive wear at the diametrical corners formed between the land clearance surface or margins on the flutes and the drill point cutting edges. Once wear begins at these corners it propagates rapidly and results in a drill that is rendered unusuable. Furthermore, conventional fluted twist drills produce the cutting action only at the drill point cutting edges and at the chisel edge. It is known that conventional twist drills are end-cutting tools in which the flutes are employed for the admission of cutting fluid and coolant and for the ejection of the chips. Normally, conventional twist drills divide into heavy duty and high helix or fast spiral drills. Heavy duty drills have great strength in the web and possess great rigidity which is suitable for severe drilling conditions in tough alloy steel, steel forgings and hard cast ferrous materials. The high-helix drills have a higher helix angle and are intended for deep holes where clearing chips from the hole is important.\nHeat is the primary cause of drill failure. The major place for heat to be generated is at the intersection of the cutting edges and the margins. When a drill gets hot it softens the material and a plating action occurs between the material being drilled and the drill. As the plating progresses, friction increases since the margins have the highest peripheral speed and absorb heat rapidly. Heat introduces the further problem of thermal expansion of the drill diameter across the margins. Drilling a material like titanium results in the drill expanding faster than titanium which causes high side forces at the margins, and galling or plating occurs rapidly. If the drill has a heavy or thick core, chip clearance is restricted and the result is additional heat generation, higher torque and early drill failure. A drill sharpened off-center can encounter some or all of the problems noted above. After heat starts to rise in the drill point, a chain reaction occurs through expansion, increased friction, cutting edge softening, and plating, and each reaction fortifies the others."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nDifferent types of antennas are used in the automotive industry, including AM/FM radio antennas, satellite digital audio radio service antenna (SDARS), cellular phone antennas, satellite navigation antennas, etc.\nAutomotive antennas may be mounted inside or outside of a vehicle. For example, a satellite navigation antenna may be mounted on an exterior vehicle surface, such as the roof, trunk, or hood of the vehicle to help ensure that the antenna has unobstructed views overhead or toward the zenith. As another example, a satellite navigation antenna may be mounted inside an instrument panel of the vehicle. The satellite navigation antenna may be connected to one or more electronic devices (e.g., an in-dash touchscreen display, etc.) inside the passenger compartment of the vehicle.\nFIG. 1 illustrates a conventional antenna assembly 100 including a patch antenna 104, a printed circuit board assembly (PCBA) 108, and an electromagnetic interference (EMI) shield 112. As shown, the conventional patch antenna 104 includes a radiating element, antenna structure, or radiator 116 on a top surface of a first dielectric substrate 120. The first dielectric ceramic substrate 120 is positioned to be mounted on top of a second dielectric substrate or board of the PCBA 108.\nThe antenna assembly 100 also includes a connector 124 for electrically connecting the PCBA 108 to a communication link, which, in turn, may be connected to an electronic device (e.g., an in-dash touchscreen display, etc.) inside the passenger compartment of a vehicle. The antenna assembly 100 includes a cover or housing 128 that may be coupled to (e.g., snapped together with, latched to, etc.) a base 132. Bumpers 136 are positioned between the PCB 108 and the cover 128.\nThe patch antenna 104 may be electrically connected or coupled to the PCBA 108 via a connector 118. The connector 118 may comprise an uninsulated pin or an insulated pin (e.g., a metal conductor with an EMI shield around it, etc.). The pin may be a semi-rigid pin and extend from the top of the patch antenna 104 through the patch's radiating element 116 and first dielectric substrate 120 to the PCBA 108.\nA label 190 may be adhesively attached to an outer surface of the cover 128. The label 190 may include information for identifying the particular antenna assembly.\nFIG. 2 illustrates another conventional antenna assembly 200 including a patch antenna 204 and a PCBA 208. As shown in FIG. 2, the conventional patch antenna 204 includes a radiating element 216 on a top surface of a first dielectric ceramic substrate 220. The first dielectric ceramic substrate 220 is mounted on top of a second dielectric substrate or board of the PCBA 208. The antenna assembly 200 also includes a connector 224 for electrically connecting the PCBA 208 to a communication link, which, in turn, may be connected to an electronic device (e.g., an in-dash touchscreen display, etc.) inside the passenger compartment of a vehicle.\nThe patch antenna 204 may be electrically connected or coupled to the PCBA 208 via a connector 218. The connector 218 may comprise an uninsulated pin or an insulated pin (e.g., a metal conductor with an EMI shield around it, etc.). The pin may be a semi-rigid pin and extend from the top of the patch antenna 204 through the patch's radiating element 216 and first dielectric substrate 220 to the PCBA 208."}
{"text": "1. Field of the Invention\nThis invention relates to hydraulic fluid compositions, particularly brake fluid compositions for use in hydraulic brake systems of automobiles.\n2. Description of the Prior Art\nRecently, automobiles have tended to become faster and larger, at the same time that greater safety is required. For this purpose, hydraulic fluids of higher performance are strongly demanded.\nThe first requirement for hydraulic fluids is to be free from the so called \"vapor lock phenomenon\". This phenomenon is caused by the vaporization of hydraulic fluids and makes brake control impossible. Consequently, brake fluids having a higher boiling point are demanded. Efforts have been made to develop hydraulic fluids having a high boiling point even in the moist state and which can maintain the higher boiling point for a long period of time. The conventional hydraulic fluids which contain a high molecular weight polyether as base polymer and a low molecular weight glycol ether as diluent, are hygroscopic and tend to suffer a severe drop in their boiling points attendant upon moisture absorption. Such hydraulic fluids are therefore unlikely to pass the standard of DOT 4 (higher than 155.degree. C.) with respect to the wet equilibrium reflux boiling point (boiling point in a moist state), according to the hydraulic fluid specification of U.S. Department of Transportation [DOT].\nHitherto, there have been proposed several hydraulic fluids which contain borate esters of glycol ethers. Such hydraulic fluids may be adequate regarding their wet equilibrium reflux boiling points, and some of them have high enough wet equilibrium reflux boiling points to pass the standard of DOT 4. But these fluids have the drawback that they cause corrosion of metals.\nIt has been proposed heretofore to add to such fluids corrosion inhibitors such as alkanolamines (mono-, di- and triethanolamine and the like). By using such corrosion inhibitors prevention of corrosion for a short time may be attained; but such known corrosion inhibitors are not effective in inhibiting the metal corrosion for a long period of time. A need exists therefore, for hydraulic fluids having improved non-corrosive properties to metals for a long period of time (test such as 2000 hours at 100.degree. C.)."}
{"text": "The International Telecommunications Union (ITU), in its International Radio Regulations, divides the world into three ITU regions for the purposes of managing the global radio spectrum. Each region has its own set of frequency allocations (i.e., frequency bands). Thus, operation in a frequency band that is allowed in one region may be prohibited in another region. For example, operation in the 915 MHz (Megahertz) band (902 MHz to 928 MHz) that is permitted in ITU region 2 (Americas) under FCC (Federal Communications Commission) part 15 rules, is prohibited in ITU Region 1 (Europe, Middle East and Africa (EMEA)). The 915 MHz frequency band may be utilized for mesh networks and relatively low-cost sensor deployments, for example.\nIn addition to the frequency bands, allowable transceiver operating parameters may also vary by region. For example, a transceiver located inside a building may be allowed a higher transmitter power than a transceiver located outside the building. The rules are complex, vary by region, country and/or subregion, and may change over time. Thus, it is difficult for a designer to ensure that a device, e.g., transceiver, will comply with the regulatory rules at its deployment location(s). This problem is compounded for the Internet of Things (IoT) and/or mesh networks that include a relatively high number of devices.\nAlthough the following Detailed Description will proceed with reference being made to illustrative embodiments, many alternatives, modifications, and variations thereof will be apparent to those skilled in the art."}
{"text": "With rapid growth of mobile subscribers and incessant development of communication techniques, the unit capacity of single HLR in mobile communication system has reached tens of thousands or even more. Since HLRs have stored subscription information of mobile subscribers, mobile services for all the subscribers will cease when the HLRs fail for a long period of time due to unforeseeable causes such as power failure, fire, earthquake or lightning strike, etc. Therefore, it is urgent to provide remote redundancy for the HLRs.\nThe redundancy for HLRs includes two aspects, one aspect is data redundancy simply for backup subscribers' data, the other is service redundancy which is based on the former aspect, i.e., when a master HLR breaks down, a redundant HLR is able to take over the services by means of signaling switchover according to the backup subscribers' data, so as to prevent the calling services from being interrupted at a maximum level.\nThe topology structures of data redundancy for HLRs are classified into two types. One is 1+1 redundancy solution, and the other is N+1 redundancy solution.\nThe so-called 1+1 redundancy solution is to provide a redundant HLR for each master HLR in the existing network so as to achieve real-time backup for external services. Though the 1+1 redundancy solution is able to take over the services seamlessly by adopting the redundant HLR when the master HLR fails, apart from some extremely important places, the application field of this solution is relatively limited in consequence of its high investment and low rate of utilization.\nThe N+1 redundancy solution is to set up a special redundant center HLR for providing redundancy for N master HLRs. In the N+1 redundancy solution, the redundancy can be further divided into compatible redundancy and non-compatible redundancy according to the result whether the master HLRs and the redundant center HLR belong to the same manufacturer. Because there are many HLR equipment providers, and different HLRs have different data saving formats and data processing methods, the compatible solution is more widely used. Since the N+1 compatible redundancy solution is unable to back up subscribers' dynamic data in real time, it will lead to the situation that the subscribers cannot act as the called party in a period of time after a failure has appeared and the services have been taken over. However, as one redundant center is able to back up data for a plurality of master HLRs made by different manufacturers, the return rate on investment of the N+1 compatible redundancy solution is relatively better, thus this solution is broadly adapted in the network.\nAt present, the N+1 compatible redundancy solution mainly has two types, one is a baseline synchronization solution, and the other is a BOSS (Business Operation Support System) increased command synchronization solution.\nWhat is called the baseline synchronization solution is to output the data files containing all subscribers' subscription information, i.e., the service data baseline of subscriber, by the master HLRs periodically, then convert the format of these files automatically by the redundant center into the format recognizable by the redundant center HLR, and load the files into the redundant center HLR. Thereby, it will realize the synchronization of subscribers' data between the master HLRs and the redundant center HLR. As the baseline synchronization solution requires the master HLRs to output the service data baseline of all subscribers periodically, performance of the host machine of the redundant center HLR is affected greatly, and it takes quite a long time to perform conversion and synchronization for the subscribers' data, it is seldom used in practical engineering.\nThe BOSS increased command synchronization solution is to convert the successfully executed business operation commands which were formerly sent to the master HLRs for modifying the subscribers' subscription data according to the interface format of the redundant center HLR, and load the commands into the host machine in real time via the business operation interface of the host machine of the redundant center HLR. Thus, as long as the master HLRs and the redundant center HLR implement baseline synchronization for one time in initial state, the conformity of the subscribers' data between the master HLRs and the redundant center HLR will be guaranteed by means of this way to achieve the synchronization of the subsequent increased command or data.\nSince the BOSS solution does not have any special requirements on the performance and function of the host machine of the redundant center HLR, and only needs to add converting function for operation commands of different manufacturers' formats in the BOSS, this solution is simple and is low cost because of low engineering work required. Therefore, the solution is widely used in the network and has gradually turned into a mainstream of N+1 redundancy solution.\nHowever, only operation commands processed by the business hall can be synchronized automatically by the BOSS solution, a large number of other operation commands processed out of the business hall, such as the operation commands for setting a batch of subscriber services in a Local Maintenance Terminal (LMT) of HLR, and setting parameters of HLR in the system level, can merely be synchronized manually by the administrator in terms of the procedure. Therefore, the maintenance workload of this solution is heavy."}
{"text": "Copper has higher conductivity and better electromigration resistance property than aluminum. Copper is thus widely used as interconnecting wires in Very Large Scale Integration (VLSI) devices. However, copper tends to diffuse in a dielectric layer and causes leakage current and a breakdown of the dielectric layer. Therefore, a barrier layer is disposed between copper interconnecting wires and the dielectric layer to prevent the diffusion of copper. With the development of VLSI devices, especially with the scaling down of high-performance logic devices, the diffusion of copper into the dielectric layer that is disposed between adjacent interconnecting wires can cause a breakdown of the dielectric layer.\nAccording to the breakdown characteristics of the dielectric layer, the breakdown can be divided into two types: an intrinsic breakdown and a Time Dependent Dielectric Breakdown (TDDB). When a voltage is applied to two adjacent copper interconnecting wires that are separated by a dielectric layer, an electric field is generated in the dielectric layer disposed between the two adjacent wires. If the generated electric field has a strength equal to or greater than a critical field strength of the dielectric layer in a copper interconnecting structure, a current will flow in the dielectric layer, and an intrinsic breakdown may occur. However, when the electric field strength is less than an intrinsic breakdown field strength, the intrinsic breakdown does not immediately occur, but the time dependent dielectric breakdown will be induced over a certain period of time due to a weakness in the isolation in the dielectric layer over time under the bias condition.\nFurthermore, the many causes for the time dependent dielectric breakdown are described as follows. As the integration degree of chips increases, interconnecting wires become so thin that current density in the interconnecting wires increases exponentially. Under the influence of the high current density, metal ions in the interconnecting wires migrate in a direction opposite to a direction of the electron motion. This phenomenon is called electromigration. In the electromigration phenomenon, migration of the metal ions may induce pileup in a local area so that hillocks are formed in a metal layer of the interconnecting wires. The electromigration may also induce voids in the interconnect wires because of mass loss, thereby degrading the interconnect performance and causing opens to the interconnecting wires.\nTherefore, preventing the loss of copper ions in the metal layer in copper interconnecting wires can improve the time dependent dielectric breakdown. In a process for forming the copper interconnecting wires, the copper interconnecting wires are exposed in an environment including an etching gas having an oxidizing ability so that copper atoms on the surface of the copper interconnecting wires are oxidized to CuO. To solve the problem mentioned above, a method for deoxygenizing copper atoms using N2 or H2 is disclosed in a paper entitled “Avoiding Cu Hillocks during the Plasma Process” published in Journal of The Electrochemical Society co-authored by Tsung-Kuei and Kang. The method for deoxygenizing copper atoms by using N2 or H2 is based on the principle: plasma is ionized to be ions or atoms in high voltage environment, then the ions and atoms can deoxygenize Cu from CuO in a deoxidization reaction on the surface of the copper interconnecting wires. However, the copper atoms deoxygenized are still in an unstable state so that the efficacy of improving the time dependent dielectric breakdown by inhibiting the loss of copper ions is not evident. Furthermore, in a plasma etching process in which a trench is formed in an upper interconnecting layer, the plasma may react with a second dielectric layer on the upper interconnecting layer, in which defect appears on a surface of the second dielectric layer. Although the second dielectric layer generally includes low dielectric constant materials, the defect will make a dielectric constant of the second dielectric layer become greater in a subsequent process of forming the copper interconnecting wires.\nIn view of analysis above, there is a need to provide a post-etching treatment process for the copper interconnecting wire, which can prevent the loss of the copper ions in the metal layer in the copper interconnecting wires, thereby avoiding the time dependent dielectric breakdown and reducing the defect on the surface of the second dielectric layer."}
{"text": "1. Field of the Invention\nThe present invention relates to a tension arm driving apparatus for a video cassette tape recorder, and particularly to an improved tension arm driving apparatus for a video cassette tape recorder capable of achieving a stable running state of a tape by preventing a tape jog in a high speed tape rewinding mode.\n2. Description of the Conventional Art\nReferring to FIG. 1, the conventional tension arm driving apparatus for a video cassette tape recorder includes a tension arm 2 rotatably engaged to a shaft pin 1 fixed to a predetermined portion of a main base 20, a tension post 3 fixed to an end of the tension arm 2, a tension spring 4 engaged to a predetermined portion of the tension arm 2 for elastically supporting the tension arm 2 in a certain direction, a cam gear 7 rotatably engaged to the shaft pin 6 and spaced apart from the tension arm 2 in cooperation with a cam groove 5, and an operation plate 10 linearly reciprocating above the main base 20 in cooperation with the cam pin engaged to an engaging pin 8 contacting with one end 2a of the tension arm 2 and the cam groove 5.\nThe operation of the conventional tension arm driving apparatus for a video cassette tape recorder will now be explained.\nTo begin with, when the cam gear 7 rotates, since the cam pin 9 of the operation plate 10 is inserted into the cam groove 5 formed on a predetermined portion of the cam gear 7, the operation plate 10, as shown in FIG. 1, moves either in the rightward or leftward direction and drives the tension arm 2.\nFor example, when the operation plate 10 moves in the leftward direction, as indicated by the arrow as shown in FIG. 1, since the end 2a of the tension and 2 engages with the engaging pin 8 of the operation plate 10, the tension arm 2 is moved by the elastic force of the tension spring 4 and rotates in the clockwise direction about the shaft pin 1.\nHowever, on the contrary, when the operation plate 10 moves in the rightward direction, the tension arm 2 rotates counterclockwise about the shaft pin 1 by the recovering force of the tension spring 4 about the shaft pin 1.\nThe tension post 3 fixed to the tension arm 2 guides the running operation of a tape T in accordance with the rotation operation of the tension arm 2. In case of using a 8 mm cameorder tape damage occurs due to shape of the cassette while the tape T is conveyed therein without a certain loading therein, due to the shape of the cassette. So, as shown in FIG. 1, the tape T is conveyed in a state where the tape T is presented therein.\nHowever, the conventional tension arm driving apparatus for a video cassette tape recorder has disadvantages in that since the tension arm moves in only one direction by the engaging pin of the operation plate, the tension arm becomes jogged in the tape rewinding mode by the tension force of a running tape, so that the running state of the tape becomes unstable. In particular, in case of using a 8 mm tape, malfunction may occur in the system because the tape comes off from its running path."}
{"text": "The dissemination of information to multiple users can be classified into two general categories. In one category, the information is managed at a central facility that is accessed by the users when they require the information. The information may be resident on multiple information servers that are geographically remote from one another, but are accessed through a central point by the users. In this type of system, the information is not replicated, except when a user manually saves a document in local storage, or the data is temporarily stored in cache memory as part of the information retrieval mechanism. In this type of system, only one master copy of the information exists, and therefore control over its content is easily maintained. An example of this type of information dissemination system is an Internet website.\nAn inherent requirement of this category of system is that the user must be connected to the central facility in order to access the information. As a result, a disconnected user, e.g. an employee traveling on an airplane, is limited in terms of the ability to obtain the information when needed. Even when a user is connected, excessive load on a server or network performance issues, such as low bandwidth or high latency, can result in inefficient access. Accordingly, the second general type of information dissemination system employs data replication. In this type of system, copies of the information are disseminated to, and stored at, multiple nodes. Once the copy has been stored at a node, the user is no longer required to be connected to the information source in order to access the information. Although advantageous from this standpoint, replication systems present a number of other issues that must be addressed. Foremost among these is the need to maintain consistency among the replicated copies of the information at the nodes. If changes are made to the information, these changes need to be disseminated to all of the nodes. Accordingly, replication systems employ different approaches for synchronizing the information among the multiple nodes.\nOne approach employs a single master version of the information, which is disseminated from a central source to replicas on remote devices that form the local nodes. The data flow is in one direction only, from the central master to the remote devices. In this type of system, all write operations occur at the master, and the users at the remote devices have read-only privileges. The set of information that is stored on each remote device is regarded as being no larger and no older than the central master.\nWhen the number of content recipients is large, the central source can become overloaded. Furthermore, there is the possibility that a large number of nodes might be connected to the central source by a common network that can experience bandwidth bottlenecks. To address these types of performance issues, hyperdistribution systems might be desirable. In this type of arrangement, the information content is disseminated from the central source to a plurality of intermediate distribution points, and from there to the ultimate recipients at the remote locations.\nAlternatively, a multiple master approach may be desirable, where users at the remote devices have the ability to modify the data. In this approach the potential exists for conflicting versions of the data at different nodes. Some multi-master systems address this issue by utilizing conflict avoidance, whereas others employ conflict resolution. To achieve conflict avoidance, a single user or application must be able to lock access to a data object before a write operation takes place, so that plural simultaneous edits cannot occur. To obtain a lock, however, connection to a resource which manages the locks is necessary. Hence, conflict avoidance is most practical in systems where the modifications are relatively infrequent.\nIn contrast, conflict resolution systems permit multiple simultaneous edits to occur, and conflicts are reconciled afterward. One type of dissemination system that is capable of resolving conflicts employs hub-and-spoke synchronization. This type of synchronization is commonly found, for example, in personal information management systems, in which mobile users are able to reconcile calendar and contact information with a single point of control. In this arrangement, each user device, e.g. personal digital assistant, carries a subset of the data in a central reference repository. When a device is connected to the central repository, its data set is synchronized with the reference data set. An example of this type of system is disclosed in U.S. Pat. No. 6,295,541. In these types of systems, the central repository functions as a clearinghouse for conflict resolution. Typically, immediate resolution is provided during synchronization, for example on the basis of time stamps. The business logic associated with this approach assumes that the individual who is performing the synchronization is the best person to resolve conflicts that arise during the synchronization. The conflict resolution process involves replacement of the older version of the information with the newer version, i.e., the prior version ceases to exist after conflict resolution.\nAn alternative to the hub-and-spoke arrangement is a peer-to-peer network, in which no single device functions as a central master. In this system, any device can synchronize with any other device at any time. This type of arrangement may be more suited to operations such as team collaboration. However, conflicting data entries can appear and be resolved many times, as disparate versions of the information propagate independently through the network and collide at various locations. Unlike hub-and-spoke mechanisms, there is no single valid version of the information, with respect to which the data set on a given device can be considered to be merely out of date.\nVarious combinations of these systems might also exist. For example, a project that employs hub-and-spoke synchronization for collaboration might also require hyperdistribution from a publisher. As a further extension of this arrangement, an object from a publisher might need to be identified with a corresponding object on the user's devices, and a workflow mechanism might allow each user to accept or postpone acceptance of any revisions to that object. As another example, a handheld device might be slaved via hub-and-spoke synchronization to a laptop, that also participates in peer-to-peer synchronization with other laptops.\nMany data replication systems are based on a model that assumes that there is only one correct version of a data object at any given time. To assure a high degree of consistency among the copies of that data object, there is an implicit requirement that the nodes connect to one another on a frequent basis, so that the amount of revision to a data object between updates is relatively small. Thus, in a hub-and-spoke system, for example, since all information is expected to flow through the hub, frequent connections along the links that constitute the spokes are formed.\nWhen there is no central hub for conflict resolution, an order-dependent approach can be employed to resolve the conflicts, but this can lead to unpredictable results. For instance, if node A is synchronized with node B, and then node B is synchronized with node C, one version of the data may be generated, whereas if node A is synchronized with node C and then node B is synchronized with node C, a different version may result.\nIt may be preferable to utilize an order-independent approach that allows business logic rules to be enforced in a way that provides more consistent results. If resolution is to be accomplished immediately, however, the business logic rules may have to be somewhat limited and/or restrictive.\nIt is an objective of the present invention to provide a framework for data replication that employs a different model to support relatively long periods of disconnection, and yet is compatible with each of the various dataflow techniques for replication described previously."}
{"text": "1. Field of the Invention\nThe present invention relates generally to control systems, and more particularly, the present invention relates to adjustment and calibration of electromagnetic devices.\n2. Description of Related Art\nElectromagnetic devices are well known. One example is an E-I core actuator, which is a type of electromagnetic linear motor so named because of its two main components. The first component is the E-core, which is a three-barrel structure having a shape that resembles the letter “E” with an insulated electric coil wire wound around the center bar and a source of current supplying current to the coil. Current running through the coil creates an electromagnetic field which attracts an associated I shaped core. Thus, an electromagnetic force is exerted across the width of a gap between the E-core and the I-core. When a constant current is supplied to the coil, the force of the electromagnetic field may change as the gap distance changes. This change in force is often referred to as the output force gain of the E-I core system.\nE-I core electromagnetic devices may be used to precisely adjust the position of an object. Unlike, for instance, a bi-directional voice coil motor which also provides precision positioning, E-I core electromagnetic devices use substantially less electric current and therefore less energy in the form of waste heat. Another benefit of E-I core electromagnetic devices is the reduction of vibration during precision motion. For instance, precision motion is frequently needed in machining, lithography, and other strict tolerance manufacturing applications e.g., in stepper and scanner machines used in the semiconductor industry. Typically, the goal is to provide precise adjustment of, for instance, a sampler or work piece stage in three dimensions.\nIn the prior art, calibration and adjustment is often done through a mechanical adjustment, which has been found to be time consuming and imprecise, especially due to problems of drift attributable to thermal or other effects. This both degrades performance and reduces system throughput, since time is required for the actual calibration. An improved calibration method would be very desirable for such systems.\nWhen using an E-I core electromagnetic device, the output force gain may be used to calibrate the device for precision adjustments, such as for positioning components of a precision machine. Precisely positioning machine components is difficult because the output force gain varies due to effects such as part-to-part variance, geometric mounting inaccuracy, and a dynamically changing gap distance during operation of E-I core electromagnetic device. Using a gap distance measurement, E-I core commutation equations may be used to model the output force gain. These models require burdensome hand tuning for each E-I core electromagnetic device. Further, even with a model, it is difficult to obtain precise measurements when dealing with large ranges of gap distance. Moreover, equations used in models, such as force gain model equations, become ineffective when gap distance information is unavailable or the force gain changes due to unmodeled factors.\nThus, there is a need for an improved method of modeling output force gain in electromagnetic devices to create precise measurements and adjustments. Further, there is a need for an improved method of calibrating the output force gain for E-I core electromagnetic devices."}
{"text": "Packaging and displaying a product can be an important way for the producer/manufacturer to feature the product and to distinguish it from the competition. For example, a packaging arrangement which permits extensive display and graphics capabilities sometimes plays a significant role from a marketing perspective. It is therefore desirable to provide a package which allows the package to be easily displayed and to have a large enough surface area to allow for graphics.\nOne way of providing such display capabilities is to include a hanghole in the top of the package to permit the arrangement to be hung from a hook for vertical display.\nIt is also advantageous to provide features in the packaging which consumers will appreciate for convenience and safety. For example, it is often desirable for many consumers to be able to easily and effectively reseal a package after it has been opened. One way of accomplishing that is through reclosable zipper seals. Reclosable seals alone, however, do not provide a way to indicate that the package may have been tampered with before the purchase. For this reason, packages with reclosable zippers have been permanently sealed above or below the reclosable zipper in such a manner that the opening of the package becomes apparent.\nIn the manufacture and use of sealed packages, there are a number of drawbacks. One drawback is that while the seal inhibits tampering, it also makes it difficult for the consumer to open the package or break the seal. Another drawback is the competition between the large surface area to display graphics and the need for adequate material for the hanghole. Further, there should be sufficient area above the hanghole to prevent the weight of the package from tearing through the hanghole. Yet another drawback in the design of such packages is the need to contain costs by minimizing the amount of film used to manufacture the package.\nAccordingly, there is a need for a packaging arrangement which overcomes the aforementioned drawbacks and yet realizes production-related cost objectives."}
{"text": "The invention relates to a temperature measuring circuit comprising a plurality of temperature sensors, each of which is formed by a temperature-dependent measuring resistance having two current terminals and two voltage tapping terminals each arranged between an associated current terminal and the temperature-dependent measuring resistance, and an evaluating circuit which via the current terminals sends a constant current through each temperature sensor and measures the voltage drop caused by the constant current between two voltage conductors connected to the voltage tapping terminals.\nIn temperature measuring circuits of this type the temperature measurement is based on the measurement of the resistance value of the temperature-dependent measuring resistance, the resistance measurement taking place by current-voltage measurement in a four-wire circuit. The use of the four-wire technique gives a higher measuring accuracy because voltage drops caused by the current flowing through the measuring resistance at lead and contact resistances do not enter the measurement result. A disadvantage is that each temperature sensor is connected by four conductors to the evaluating circuit and consequently with a relatively large number of temperature sensors, as are often required for detecting different temperatures at different points in technical plants or processes, the number of connection conductors required rises in proportion with the number of temperature sensors and often becomes undesirably large.\nThe problem underlying the invention is to provide a temperature measurement circuit of the type set forth at the beginning which whilst completely retaining the advantages of the four-wire technique requires a substantially smaller number of connecting conductors between the temperature sensors and the evaluating circuit.\nAccording to the invention this problem is solved in that the temperature sensors are connected via the current terminals in series, that the two voltage tapping terminals of two temperature sensors which follow each other in the series circuit and are associated with two current terminals connected together are connected to a common voltage conductor, and that a resistance is inserted into the connection between at least one of the two voltage tapping terminals and the common voltage conductor.\nIn the temperature measuring circuit according to the invention due to the series connection of the temperature sensors only two current conductors are required for the connection to the current source, irrespective of the number of said sensors, and the number of the voltage conductors required is reduced in that in each case two voltage tapping terminals belonging to two consecutive temperature sensors in the series circuit are connected to a common voltage conductor instead of to two separate voltage conductors. The total number of connecting conductors required is therefore only three greater than the number of the temperature sensors compared with four times the number of temperature sensors in conventional temperature measuring circuits. Accordingly, the saving of connecting conductors is the greater the greater the number of temperature sensors present. Due to the fact that each common voltage conductor is connected to the two associated voltage tapping terminals not directly but via resistances, it is ensured that the voltage measurement takes place practically with the same accuracy as with the conventional four-wire technique with separate voltage conductors.\nAdvantageous further developments and embodiments of the invention are characterized in the subsidiary claims."}
{"text": "Many plasma diagnostic techniques use a swept-wavelength laser to obtain atomic absorption spectra from a plasma stream. Most of these measurements are dependent upon the scan width of the laser, typically measured in GHz. Since there are no instruments that measure scan width directly, the laser operator must attempt to manually adjust the scan width to specifications (within a few percent), and keep it there at all times. Results of many spectroscopic measurements, such as plasma density, are directly proportional to scan width. Two instruments are presently available to assist the operator in this scan width adjustment effort: wavemeters and Fabry-Perot confocal interferometers (F-Ps).\nWavemeters directly measure absolute wavelength by comparison with a reference light source or with a precision grating. However, only the most expensive and exotic wavemeters have sufficient resolution to measure the difference between the endpoints of a narrow scan, and even the best cannot track the wavelength as it rapidly changes.\nAnother way to directly measure scan width employs Fabry-Perot confocal interferometers (F-Ps). The F-P samples the laser beam at some point in the optical system. As the laser is scanned, the F-P interferometer produces a voltage pulse each time the wavelength crosses through a free-spectral resonant mode of the interferometer's cavity. The laser operator can observe the F-P pulses on an oscilloscope and attempt to visually count them, then adjust the laser until the proper number appears. A simple digital counter, synchronized with the scan, could also help to automate the counting process. With a Coherent, Inc., model 216 Fabry-Perot interferometer, a pulse occurs every 300 MHz. For example, with a scan of 3 GHz, 10 F-P intervals appear during each scan. The laser operator could then trim the scan width by manually adjusting the scan width knob on the laser's control chassis. However, the precision of this measurement is only .+-.10% (one in 10). Thus, a 10% peak error in any wavelength-dependent measurements may be expected. Also, the operator must periodically adjust the laser to compensate for drift. For this reason, an improved, automated measurement and control system for scan width using a real-time control computer and some specialized CAMAC (Computer Automated Measurement and Control, IEEE 595) modules designed specifically for this application would be desirable. The present invention presents such a system."}
{"text": "The present invention relates to resealing data. More specifically, the present invention relates to updating reduction ratios used for rescaling data sequences. Still more specifically, the present invention provides techniques for determining reduction ratios for modifying transform coefficients associated with an input data sequence (e.g. an audio segment or a video sequence) to provide modified transform coefficients associated with a modified output data sequence.\nVideo data is one particularly relevant form of data that can benefit from improved techniques for resealing. Generally, compressing data or further compressing compressed data is referred to herein as resealing data. Video rescaling schemes allow digitized video frames to be represented digitally in an efficient manner. Rescaling digital video makes it practical to transmit the compressed signal by digital channels at a fraction of the bandwidth required to transmit the original signal without compression. International standards have been created on video compression schemes. The standards include MPEG-1, MPEG-2, MPEG-4, H.261, H.262, H.263, H.263+, etc. The standardized compression schemes mostly rely on several key algorithm schemes: motion compensated transform coding (for example, DCT transforms or wavelet/sub-band transforms), quantization of the transform coefficients, and variable length coding (VLC).\nThe motion compensated encoding removes the temporally redundant information inherent in video sequences. The transform coding enables orthogonal spatial frequency representation of spatial domain video signals. Quantization of the transformed coefficients reduces the number of levels required to represent a given digitized video sample and reduces bit usage in the compression output stream. The other factor contributing to rescaling is variable length coding (VLC) that represents frequently used symbols using code words. In general, the number of bits used to represent a given image determines the quality of the decoded picture. The more bits used to represent a given image, the better the image quality. The system that is used to compress digitized video sequence using the above described schemes is called an encoder or encoding system.\nMore specifically, motion compensation performs differential encoding of frames. Certain frames, such as I-frames in MPEG-2, continue to store the entire image, and are independent of other frames. Intracoded frames, such as B-frames or P-frames in MPEG-2, store motion vectors associated with the movement of particular objects. The pixel-wise difference between objects is called the error term, which can be stored in P-frames and B-frames. In MPEG-2, P-frames reference a single frame while B-frames reference two different frames. Although this allows fairly high reduction ratios, motion compensation is limited when significant changes occur between frames. This precludes high reduction ratios. Furthermore, motion compensation can be computationally expensive.\nEach frame can be converted to luminance and chrominance components. As will be appreciated by one of skill in the art, the human eye is more sensitive to the luminance than to the chrominance of an image. In MPEG-2, luminance and chrominance frames are divided into 8×8 pixel blocks. The 8×8 pixel blocks are transformed using a discrete cosine transform (DCT) and scanned to create a DCT coefficient vector. Quantization involves dividing the DCT coefficients by a scaling factor. The divided coefficients can be rounded to the nearest integer. After quantization, some of the quantized elements become zero. The many levels represented by the transform coefficients are reduced to a smaller number of levels after quantization. With fewer levels represented, more sequences of numbers are similar. For example, the sequence 4.9 4.1 2.2 1.9 after division by two and rounding becomes 2 2 1 1. As will be described below, a sequence with more similar numbers can more easily be encoded using VLC encoding. However, quantization is an irreversible process and hence introduces significant loss of information associated with the original frame or image.\nVLC encoding takes the most common long sequences of numbers of bits and replaces them with a shorter sequence of numbers or bits. Again, VLC encoding is limited by common sequences of numbers or bits. Data containing fewer common sequences take more bits to encode.\nCurrently available compression techniques for rescaling data (e.g. video or audio) are limited in their ability to effectively compress data sequences for transmission across networks or storage on computer readable media. The available techniques also have significant limitations with respect to loss, computational expense, and delay. Various techniques for reducing the bit rate of compressed data sequences including audio and video streams are being developed. Some of the more promising approaches are described in U.S. Pat. No. 6,181,711 titled System And Method For Transporting A Compressed Video And Data Bitstream Over A Communication Channel. Other approaches are described in U.S. patent application Ser. No. 09/608,128 Methods And Apparatus For Bandwidth Scalable Transmission Of Compressed Video Data Through Resolution Conversion and U.S. patent application Ser. No. 09/766,020 titled Methods For Efficient Bandwidth and U.S. patent application Ser. No. 08/985,377 titled System And Method For Spatial Temporal-Filtering For Improving Compressed Digital Video Scaling Of Compressed Video Data. Each of these references is assigned to the assignee of this invention and is incorporated herein by reference for all purposes. It is still desirable to provide additional techniques for rescaling data that improve upon the limitations of the prior art."}
{"text": "Existing multi-frequency wireless devices (e.g., radios) use an antenna structure that attempts to radiate at optimum efficiency over the entire frequency range of operation, but can really only do so over a subset of the frequencies. Due to size constraints, and aesthetic design reasons, the antenna designer is forced to compromise the performance in some of the frequency bands. An example of such a wireless device could be a mobile telephone that operates over a range of different frequencies, such as 800 MHz to 2200 MHz. The antenna will not radiate efficiently at all frequencies due to the nature of the design, and the power transfer between the antenna, the power amplifier, and the receiver in the radio will vary significantly.\nAdditionally, an antenna's performance is impacted by its operating environment. For example, multiple use cases exist for radio handsets, which include such conditions as the placement of the handset's antenna next to a user's head, or in the user's pocket or the covering of an antenna with a hand, can significantly impair wireless device efficiency. Further, many existing radios use a simple circuit composed of fixed value components that are aimed at improving the power transfer from power amplifier to antenna, or from the antenna to the receiver, but since the components used are fixed in value there is always a compromise when attempting to cover multiple frequency bands and multiple use cases."}
{"text": "Granular media filters are commonly used to remove particulate solids from water in water and wastewater treatment plants. These filters typically employ sand, crushed coal, granular activated carbon (GAC), crushed garnet or a combination of these or similar granular materials as filter media. The water to be filtered is passed through the filter media and the particulate solids are trapped at the surface of the filter media and in the interstitial spaces within the filter media. The entrapped particulate solids are removed from the filter media by a periodic washing operation called backwashing in which water or a water/air combination is passed upwardly through the filter media at sufficient velocities to carry the particulate solids upwardly. The backwash water and entrained particulate solids are intercepted by a backwash water collector typically consisting of a pipe or trough. Because of turbulence caused by the backwashing process, grains of the filter media can become entrained in the backwash water stream and are carried over the overflow weir edge of the backwash trough and are removed or lost from the filter. The media loss is particularly severe when using low density media such as crushed coal or when using air to assist in scouring the entrapped particulate solids from the filter.\nOne prior art granular media filter is shown in U.S. Pat. No. 4,076,625 and includes means intended to reduce media loss and consisting of baffles placed adjacent to or surrounding the backwash trough to intercept and deflect the air used in the backwash process. The baffles reduce the velocity and turbulence of water currents near the overflow weir edge of the backwash trough to thereby reduce media loss."}
{"text": "For example, it is known that a machine structural component usable for a power transmission part, such as a gear for car transmission, may develop pitching damage, i.e., peel failure of a metal contact portion due to increased pressure of contacted surface during use thereof. Thus, a steel component for use in such an application includes a component prepared in such a manner that one of various types of case-hardening steel such as SCr, SCM, and SNCM is subjected to forming through hot forging and cutting, and is then subjected to surface hardening treatment such as carburizing treatment or carbonitriding treatment, and furthermore a solid lubrication film such as a molybdenum disulfide film is provided on a surface of the component.\nIn recent years, however, a machine structural part is increasingly demanded to have higher output power, smaller size, and lighter weight, and thus the machine structural element for use in such a power transmission part tends to increasingly receive a higher load. As a result, it is becoming difficult to achieve the required pitching resistance not only by the component including one of various types of case-hardening steel such as SCr, SCM, and SNCM subjected to surface hardening treatment, but also by the component on which the solid lubrication film is further provided.\nIn an electric car, which is recently increasingly produced in light of reduction in environmental load, since rotation of a motor is directly transmitted to a reduction gear, such components are subjected to higher rotation than in a gasoline-powered vehicle. Furthermore, since a lubricating oil for use in the electric car has a kinematic viscosity in use environment lower than that for use in the gasoline-powered vehicle, a thin oil film is formed on a surface of each of such steel components configuring the power transmission part. In such environment, a site having almost no oil film thereon may be locally formed. In particular, along with increase in rotation and/or slippage, oil temperature increases, and kinematic viscosity of the lubricating oil tends to be lowered, and the sites each having almost no oil film thereon further increases. In such environment, therefore, steel components easily wear due to metal contact therebetween, and each steel component is easily softened through temperature rise due to friction heat; hence, seizing is likely to occur in an early stage.\nVarious technologies have been proposed on steel members usable in the environment as described above. For example, PTL 1 discloses a component having high resistance to pressure of contacted surface, which is composed of steel that contains C: 0.15 to 0.40%, Si: 0.50 to 1.50%, Mn: 0.20 to 1.50%, Cr: 0.50 to 1.50%, Mo: 0.05 to 0.50%, at least one element selected from the group consisting of Ni: 0.50 to 3.50%, Ti: 0.03 to 0.20%, Nb: 0.03 to 0.15%, and Al: 0.01 to 0.10%, and P: 0.010% or less, with the remainder consisting of Fe and inevitable impurities, wherein a rolling site of a surface of the component has a carbon content of 0.8 to 1.2%.\nHowever, in this technology, the carbon content is provided through carburizing or carbonitriding treatment without nitriding treatment. Hence, the nitrogen amount of a surface layer is small, less than 1%. Consequently, if this component is used in an environment where seizing more easily occurs, the component cannot maintain excellent seizing resistance.\nPTL 2 discloses a method of manufacturing a machine structural component having excellent fatigue strength, particularly surface fatigue strength, the method being characterized in that steel is used as a material, the steel containing C: 0.4 to 0.7%, Si: 0.3% or less, Mn: 0.2 to 1%, Cr: 0.2 to 3%, Mo: 0.1 to 1%, V: 0.1 to 1%, Al: 0.01 to 0.05%, N: 0.003 to 0.02%, S: 0.07% or less, and Ti: 0.002% or less, with the remainder consisting of Fe and inevitable impurities including P and O that are controlled, respectively, to be 0.02% or less and 0.002% or less, the material is formed into a predetermined component shape, and is then subjected to nitriding treatment or nitrocarburizing treatment and successively subjected to induction hardening treatment so that nitrogen is diffused from a surface of the component, and a surface-hardened layer having a nitrogen content of 0.05% or more is formed at a depth position of at least 0.2 mm from a top surface of the component.\nHowever, since this technology includes the diffusion treatment of surface nitrogen, the surface layer cannot maintain high nitrogen content. In addition, since carbide is not dispersed in a surface layer of the component, the component cannot exhibit good seizing resistance in a high slippage environment.\nPTL 3 discloses a method of manufacturing a rotary hook of a sewing machine, the method being characterized in that case-hardening steel such as carbon case-hardening steel and chromium-molybdenum case-hardening steel is subjected to carburizing-and-quenching and tempering, and further subjected to nitrocarburizing treatment. In this technology, however, only typical carburizing treatment is performed, and carbide does not dispersedly exist in a surface layer of a component. Hence, the component cannot exhibit excellent seizing resistance in a high slippage environment."}
{"text": "Today's corporate computing environments maintain huge amounts of data pertaining to business information. Business information may be virtually any kind of data needed to maintain a business and to comply with regulations regarding a regulated business. In the area of import and export, large amounts of data must be maintained for long periods of time. This data must remain accessible for: compliance screening, auditing, transaction tracking, and other business reasons.\nSome of the complexities associated with maintaining and using vast amounts of data stored for long periods of time include, but are not limited to, data table structures and databases changing over time, hard-coded user interface screens which are non-flexible, limited pre-defined “customization” options in a user interface, etc. User interface screens are typically designed by a software developer at a user-interface design time. User-interface design time typically occurs long before actual end-user usage. Different methods for allowing end-user customization are utilized by designers and can be intended to predict a user's future needs. However, over time, screens may become outdated and show data not of interest or may not allow required user controls over certain fields. Some possible side effects of an outdated information screen may be that the screen requires excessive panning and scrolling by an end-user or there may be limited filtering capability for data on a screen (as well as many other inconveniences).\nThese and other problems largely arise because a developer at design time (well ahead of actual end-user usage) is attempting to implement a “one-size fits all” capability. Also, the designer is typically forced to base designs on requirements of “typical” rather than individual users. Additionally, desires of an end-user performing a particular function may change over time as the responsibilities associated with their job change. In a corporate environment, requests to update a user interface are typically provided by an end-user filling out a change request and providing that change request to a development team. The end-user is then stuck with using the old and potentially cumbersome outdated interface while the development team prioritizes multiple change requests and produces a new version of the user-interface for a potentially large number of users.\nAnother example of how user interfaces to data may become obsolete can occur when the structure of the data storage changes over time. Historical data may be migrated in some form or fashion to the new structure (possibly without all fields of information) or maintained in two different data sources (e.g., old database/tables for old data and new database/tables for new data). User interfaces of the prior art may not be flexible enough to adapt to changing locations of data sources.\nIn summary, users accessing this vast amount of data may have different requirements and preferences as to how a user desires to interact with the data. Prior art techniques exist to provide limited end-user customization for applications and screens used to access and view data. Because of these concerns and other limitations a more general method and system for an end-user to customize how their particular interface functions may be desirable. Customization may include an ability to customize what data is shown on the screen and what capabilities are available on each screen to interact with that data. Also, a user may wish to customize what is shown to limit it to data of interest without necessarily being subjected to data fields not pertaining to a task at hand. Capabilities to provide these and other customization options to an end-user are disclosed below in a non-limiting embodiment of an import export trade transaction storage and retrieval system."}
{"text": "Free space optical switching systems and photonic computing systems utilize macroscopic optical elements such as holograms, gratings, lenses, mirrors, and fiber optic arrays as basic building blocks. In such systems information is carried by arrays of light beams that are collimated, manipulated, and focused onto spatial light modulators in a stage by stage fashion. Fiber optic arrays may serve as input terminals, output terminals and interstage couplers for such systems. Precision fiber optic arrays are important support components of such systems.\nExamples of known optical switching systems and photonic computing systems are known from: \"An All Optical Shift Register Using Symmetric Self Electro-optic Effect Devices,\" published in OSA Proceedings on Photonic Switching, Vol. 3, pp. 192-195, 1989, edited by J. E. Midwinter and H. S. Hinton, \"An All Optical Realization of a 2.times.1 Free Space Switching Node,\" published in Photonics Technology Letters, No. 8, 1990, pp. 600-602 by E. Kerbis, T. J. Cloonan, and F. B. McCormick; \"An All Optical Implementation of a 3D Crossover Switching Network,\" published in Photonics Technology Letters, No. 6, 1990, pp. 438-440 by T. J. Cloonan, M. J. Herron, F. A. P. Tooley, G. W. Richards, F. B. McCormick, E. Kerbis, J. L. Brubaker, and A. L. Lentine; \"Module for Optical Logic Circuits using Symmetric Self Electro-Optic Effect Devices,\" published in Applied Optics, No. 14, 1990, pp. 2164-2170 by M. E. Prise, N. C. Craft, R. E. LaMarche, M. M. Downs, S. J. Walker, L. A. D'Asaro, and L. M. F. Chirovsky; \"Parallel Interconnection of Two 64.times.32 Symmetric Self Electro-optic Effect Device Arrays,\" published in Electronics Letters, No. 20, 1991, pp. 1869-1871 by F. B. McCormick, F. A. P. Tooley, J. M. Sasian, J. L. Brubaker, A. L. Lentine, T. J. Cloonan, R. L. Morrison, S. L. Walker, and R. J. Crisci; \"Optomechanics of a Free Space Switch: The System,\" published in Optomechanics and Dimensional Stability, Proceedings of SPIE, No. 1533, 1991, pp. 97-114 by F. B. McCormick, F. A. P. Tooley, J. L. Brubaker, J. M. Sasian, T. J. Cloonan, A. L. Lentine, S. J. Hinterlong, and M. J. Herron; \"Experimental Investigation of a Free Space Optical Switching Network Using S-SEEDs,\" published in Applied Optics, Vol. 31, No. 26, 1992, pp. 5431-5446 by F. B. McCormick, F. A. P. Tooley, T. J. Cloonan, J. L. Brubaker, A. L. Lentine, R. L. Morrison, S. J. Hinterlong, M. J. Herron, S. L. Walker, and J. M. Sasian; and \"A Six-Stage Digital Free-Space Optical Switching Network Using S-SEEDs,\" published in Applied Optics, Oct. (1992) by F. B. McCormick, T. J. Cloonan, F. A. P. Tooley, A. L. Lentine, J. M. Sasian, J. L. Brubaker, R. L. Morrison, S. L. Walker, R. J. Crisci, R. A. Novotny, S. J. Hinterlong, H. S. Hinton, and E. Kerbis. These articles chronical the progress of optical and photonic system technology as it grows in complexity and functionality.\nThe potential of free space optical switching systems to interconnect a larger number of communication channels at high bit rates has spurred the development of such systems. The potential of photonic computing systems for increased processing rates and bit transfer rates through the use of optical and photonic components has spurred similar photonic computing system development. These potential benefits have motivated the development of high precision fiber optic arrays to communicate light beam borne data to and from free space optical switching systems and photonic computing systems. High precision fiber optic arrays have been easy to specify, but difficult to physically realize, especially if high precision positioning of the optical fibers is required. Several fabrication techniques have been reported in various publications. For example, Miller describes a 2-D array of optical fibers that was made by stacking a number of linear arrays of optical fibers that were supported by grooved spacers in \"A Fiber Optic Cable Connector,\" The Bell Technical Journal, Vol. 54, No. 9, 1547-1555, 1975. These spacers were manufactured by the precise etching of grooves in both sides of a silicon wafer, potting all the optical fibers in place and polishing the resulting assembly. This technique has been used to assemble a fiber array with a maximum positioning error of 10 micrometers as reported by U. Danzer, P. Kipfer, K. Zufi, J. Lindolf, and J. Schwider, in \"High Precision Two Dimensional Fibre Array in Silicon V-Groove Technique,\" Angewandte Optik, Physikalisches Institut der Universitat Erlangen, Annual Report 1992.\nIn another effort, an alignment-free assembly technique has been developed where fiber end positioning to within plus or minus 8 micrometers was achieved. This technique is described by A. Sasaki, T. Baba, and K. Iga in \"Put-in Microconnectors for Alignment-free Coupling of an Optical Fiber Array,\" IEEE Photonics Technology Letters, Vol. 4 No. 5, pp. 908-910, 1992. In this known technique, an array of sockets with centering plugs for optical fibers were micro-fabricated to achieve self-centering of each optical fiber upon insertion and to expedite assembly also.\nG. A. Koepf and B. J. Markey describe another technique involving arrays of precision holes in substrates to insert and locate optical fibers with a standard deviation of 12.6 micrometers in, \"Fabrication and Characterization of a 2-D Fiber Array,\" Applied Optics Vol. 23, No. 2, pp. 3515-3516, 1984. Fiber optical arrays using precision holes to position the optical fibers have also been described by Basavanhally in \"Opto-mechanical Alignment and Assembly of 2-D Array Components,\" Technical Digest of the IEEE Princeton Section Sarnoff Symposium, Mar. 26, 1992. G. M. Proudley, C. Stace, H. White also describe their fiber optic array fabrication technique in their article \"Fabrication of 2-D Fibre-Optic Arrays for an Optical Crossbar Switch,\" submitted to Optical Engineering.\nKoyabu, F. Ohira, T. Yamamoto, and S. Matsuo have realized a 2-D fiber optic array with a mean fiber positioning error of 3 micrometers by inserting the optical fibers into microferrules and stacking the fiber optic and microferrule assembly to create the array. This technique is described in their article \"Novel High Density Collimator Module,\" Technical Digest, Conference on Optical fiber Communication/International Conference on Integrated Optics and Optical Fiber Communication, 1993 Technical Digest Series, Vol. 4, (Optical Society of America, Washington, D.C., 1993), pp. 2-3.\nA common feature of these known techniques is that the final positioning of each optical fiber is accomplished by referencing each optical fiber to a mechanical jig or support through which the fibers pass. Thus, the precision of this mechanical jig or support limits the ultimate precision attainable. These known techniques, with their inherent mechanical limitations, are not precise enough to meet the requirements of the developing optical switching systems and the photonic computing systems. Thus, there is a need in the art to provide a high precision fiber optic array and a process for producing such an array.\nIt is an object of the present invention to provide a fiber optic array that is not limited by the precision of the mechanical support through which the optical fibers pass.\nIt is another object of the invention to provide a process for producing a high precision optical fiber array."}
{"text": "Another important feature of the invention is the use of injected compressors. Especially useful in heat pump systems according to the present invention is the type of injected compressor disclosed and claimed in the copending U.S. patent application of William H. Wilkinson, the present inventor, and James H. Saunders, Ser. No. 07/161,189, filed Feb. 26, 1988, for Crossed Piston Compressor with Vernier Offset Port Means, now U.S. Pat. No. 4,936,111, issued June 26, 1990. This copending patent relates to systems and apparatus for compressing gaseous fluids, especially refrigerant gas vapors operating in refrigeration cycles, combining reciprocating pistons in a rotating cylinder member that is mounted for rotation in a stationary frame and is driven by an external source of rotative power. The cylinder block rotates in an encircling port ring member which contains inlet, interstage, and outlet port sets. The number of port sets is greater or less than the number of cylinders, to provide a vernier effect in the timing of connections between the port sets and the cylinders.\nThe above mentioned copending patent is assigned to the assignee of the present invention. It is hereby incorporated herein by reference, and made a part hereof the same as if fully set forth herein, for purposes of indicating the background of the invention and illustrating the state of the art.\nThe term injected compression is used herein to mean the injection of high pressure vapor into the compression space of a compressor after the inlet suction of vapor from an evaporator is at least substantially complete. U.S. Pat. No. 4,332,144, Shaw, discusses the advantages of utilizing a scavenge vapor as a means of increasing the coefficient of performance of a refrigeration device and/or to increase the refrigeration (heat pumping) capacity of a given compressor displacement.\nThe copending patent cited above provides improved injection capability and also offers the option of providing two stages of compression in which the discharge of one stage of compression is fed to the second stage of compression with either or both of the stages capable of injected operation. So typical compressors according to the copending patent can be designed to have as many as four vapor inlets; the lowest pressure inlet being the output from the evaporator, the next higher pressure being the injection to the first stage of compression, the third pressure level being the interstage pressure, and the fourth being the injection to the second stage of compression.\nVarious advantages of such compressors are explained in the copending patent relative to refrigeration in general, but the patent does not specifically describe the advantages that an injected compressor of that type can have with NARB's as working fluids. The present invention includes additional novel system arrangements based on the general characteristics of that type of compressor. It also employs a novel type of control arrangement that provides a significant improvement over current heat pumps that extract heat from the ambient air for heating a dwelling. Improvement is shown with this system using a single refrigerant; and further improvement is shown for the slightly more complex system using a NARB as the working fluid.\nThe novel control arrangement for the single-stage heat pump compressor has direct application as the lower stage of a two-stage compressor system for low temperature freezer applications. One version of the system uses a single refrigerant as the working fluid; while a slightly more complex version uses NARB as the working fluid and is useful for reaching lower temperatures than the single-refrigerant system."}
{"text": "1. Field of the Invention\nThe present invention relates to a processor for processing an FM signal such as an FM stereo signal.\n2. Description of the Prior Art\nReferring to FIG. 1, there is shown a block diagram of the entire arrangement of an FM stereo radio receiver. Reference numeral 331 represents an antenna. Reference numeral 332 represents an RF amplifier. Reference numeral 333 represents an mixer for mixing an oscillation signal of a local oscillation circuit 334 and an RF signal from the RF amplifier 332 to output an IF signal. Reference numeral 335 represents an IF filter. Reference numeral 336 represents an IF amplifying circuit. Reference numeral 337 represents a limiter. Reference numeral 339 represents a 90.degree. phase shifting circuit. Reference numeral 340 represents a multiplier serving as an FM detector. Reference numeral 341 represents a multiplexer. Reference numeral 342 represents a power amplifying circuit. Reference numerals 343 and 344 represent left and right speakers.\nThe IF amplifier 336, the limiter 337, the 90.degree. phase shifting circuit 339, the multiplier 340 and the multiplexer 341 are formed in a one-chip IC device 338. To the 90.degree. phase shifting circuit 339, an resonance circuit 345 including a coil L and a capacitor C is externally attached, and to the multiplexer 341, a circuit 347 constituting a part of a voltage controlled oscillator (VCO) 346 is externally attached. The circuit includes a resistor and a capacitor.\nIn this FM radio receiver, the frequency of a received FM radio signal is converted into an intermediate frequency (IF) by the IF circuit 336, and the amplitude thereof is made constant by the limiter 337. Thereafter, the signal is transmitted through a line a directly to the multiplier 340 and is also transmitted through a line b to the multiplier 340 by way of the phase shifting circuit 330 where its phase is shifted by 90.degree.. These two signals are multiplied at the multiplier 340. The phase shifting circuit 339 shifts the phase by the capacitor C and the coil L. The VCO 346 used in the multiplexer 341 is designed to regulate a free-running oscillation frequency by a CR time constant of a circuit externally attached to the IC device 338.\nReferring to FIG. 2, there is shown another conventional FM radio receiver where ceramic resonators 350 and 351 are externally attached instead of the resonance circuit 345 and the circuit 347 of FIG. 1. In a phase shifting circuit 339 using the ceramic resonator 350, as shown in FIG. 3, three resistors R and the ceramic resonator 350 form a Wheatstone bridge, and the phase of the signal is shifted by using a reactance component of the ceramic resonator 350. The oscillation frequency of the VCO 346 depends on the ceramic resonator 351.\nHowever, in the phase shifting circuit 339 of FIG. 1, it is necessary to adjust a free-running oscillation frequency by the capacitor or the coil, and the regulation which is a mechanical regulation is liable to deterioration with age. In addition, since the coil cannot be formed in an integrated circuit (IC) device, more number of terminal pins are required and an externally-connected part is required. The VCO 346 has similar disadvantages since the free-running oscillation frequency is regulated by a CR time constant of a circuit externally attached to the IC device.\nWhen the ceramic resonators 350 and 351 are externally attached as shown in FIG. 2, not only an externally-attached part is required but also the part (ceramic resonator) is not available in some countries since it is a special part. In addition, the ceramic resonator is inferior to the former method in property, and difficult to match with an IC device. Further, since the free-running frequency of the VCO 346 depends on only the property of the resonator 351, to obtain a comparatively low frequency such as 19 kHz and 38 kHz necessary to process a composite signal, a number of frequency dividers are required for dividing (into 24 and into 12) 456 kHz of the ceramic resonator 351. This increases the chip area of the IC device."}
{"text": "Parallel processing algorithms take advantage of multi-processing computer architectures by distributing portions of the data among a plurality of processing elements. Algorithmic processing takes place until the system state reaches a point in which the necessary data becomes non-local to the processing elements. Hence, a re-distribution of data must be carried out prior to the next stage of algorithmic computation. Depending on the required state of data re-distribution, this operation can be very challenging to the communication system that ties the processing elements together.\nData re-distribution in system memory can pose additional challenges. System memory access is typically defined for a fixed set of increments and optimized for particular access patterns. As an example, a memory interface may only permit accesses of 64 or 128 bytes, and though system memory permits random access, it is understood by those skilled in the art that accessing contiguous address space is most efficient. Depending on the distribution of data within the multi-processor, and the required staging for the next phase of computation, a given processing element may require a large number of non-contiguous data accesses which are smaller than the allowable memory access sizes.\nCurrent multi-processor performance suffers under the scenario where a small amount of data is scattered among several processing or storage elements. An example of a challenging data re-distribution occurs when a large number of unique data values must be distributed among a plurality of processing or memory elements for real time processing. A reverse situation occurs when many processing or memory elements must each individually send a small amount of data to a single location on the multi-processor."}
{"text": "This invention generally relates to sleeves, and, more particularly, to sleeves used to wrap flower pots containing floral groupings and/or mediums containing floral groupings, and methods of using same."}
{"text": "Field of the Invention\nThe present invention relates generally to a crop cart and more specifically it relates to a crop cart unloading system for efficiently unloading a harvested crop from a crop cart.\nDescription of the Related Art\nAny discussion of the related art throughout the specification should in no way be considered as an admission that such related art is widely known or forms part of common general knowledge in the field.\nCrop carts are used to receive a crop from a harvester directly in the field during harvest. Crop carts can also be used to transport seed to a planter. Crop carts are typically used to receive, transport and unload various types of field crops such as, but not limited to, sugar beets, carrots, potatoes, onions, tomatoes, cucumbers, corn silage and the like. Conventional crop carts are comprised of a trailer that is pulled by a tractor through the field next to the harvester. Conventional crop carts include an open topped hopper that receives various volumes of harvested crop from the harvester in real-time. When the hopper of the crop shuttle is full, the crop cart is transported to a location in the field to unload to a truck or other transport device using an unloading conveyor at the front of the hopper. The hopper includes a floor conveyor that moves the crop forwardly towards the unloading conveyor where the unloading conveyor is angled upwardly to lift the crop upwardly from the crop shuttle into the transport vehicle.\nOne problem with conventional crop carts is that the unloading conveyor extends upwardly at a relatively steep angle thereby resulting in the crop falling downwardly on the unloading conveyor during the unloading procedure. When the crop falls downwardly on the unloading conveyor, crops such as sugar beets can become damaged. Another problem with conventional crop carts is that the front unloading conveyor extends upwardly above the upper edge of the container when in the storage position thereby creating potential problems of hitting the harvester during loading of the crop. Another problem with conventional crop carts is that they require a significant volume of pressurized hydraulic fluid from the tractor pulling the crop cart to unload the crop which some older tractors are unable to handle.\nBecause of the inherent problems with the related art, there is a need for a new and improved crop cart unloading system for efficiently unloading a harvested crop from a crop cart."}
{"text": "1. Field\nA cooling apparatus for a machine room of a refrigerator using a nacelle-shaped duct is disclosed herein.\n2. Background\nCooling apparatuses are known. However, they suffer from various disadvantages."}
{"text": "This invention relates to the fields of drilling and producing hydrocarbon wells, and to the measuring of downhole formation characteristics, and to drill string telemetry for bidirectional communication of measurement and control information between dowhhole and surface equipment, and to a surface communication system for bidirectional communication between drill string telemetry and a surface processor.\nThe advent of measurement while drilling (MWD) and logging while drilling (LWD), as well as development of surface control of special drilling processes, such as directional drilling, have been important advances in the art of drilling and producing hydrocarbon wells. These processes require communication, in both directions, between the surface and the downhole measuring and drilling equipment. At present, mud pulse telemetry is the only technique in widespread commercial use for communication while drilling, between downhole equipment and the surface. [Unless otherwise indicated, references, throughout, to “while drilling,” or the like, are intended to mean that the drill string is in the borehole or partially in the borehole as part of an overall drilling operation including drilling, pausing, and or tripping, and not necessarily that a drill bit is rotating.] In mud pulse telemetry, data is transmitted as pressure pulses in the drilling fluid. However, mud pulse telemetry has well known limitations, including relatively slow communication, low data rates, and marginal reliability. Current mud pulse technology is capable of sending MWD/LWD data at only about 12 bits per second. In many cases, this rate is insufficient to send all the data that is gathered by an LWD tool string, or is limiting on the configuration of a desired tool string. Also, mud pulse technology does not work well in extended reach boreholes. Signaling from uphole to downhole, by regulating mud pump flow, in order to control processes such as directional drilling and tool functions, is also slow, and has a very low information rate. Also, under certain circumstances, for example underbalanced drilling employing gases or foamed drilling fluid, current mud pulse telemetry cannot function.\nThere have been various attempts over the years to develop alternatives to mud pulse telemetry that are faster, have higher data rates, and do not require the presence of a particular type of drilling fluid. For example, acoustic telemetry has been proposed, which transmits acoustic waves through the drill string. Data rates are estimated to be about an order of magnitude higher than mud pulse telemetry, but still limiting, and noise is a problem. Acoustic telemetry has not yet become commercially available. Another example is electromagnetic telemetry through the earth. This technique is considered to have limited range, depends on characteristics, especially resistivity, of the formations surrounding the borehole, and also has limited data rates.\nThe placement of wires in drill pipes for carrying signals has long been proposed. Some early approaches to a wired drill string are disclosed in: U.S. Pat. No. 4,126,848, U.S. Pat. No. 3,957,118 and U.S. Pat. No. 3,807,502, and the publication “Four Different Systems Used for MWD,” W. J. McDonald, The Oil and Gas Journal, pages 115-124, Apr. 3, 1978.\nThe idea of using inductive couplers, such as at the pipe joints, has also been proposed. The following disclose use of inductive couplers in a drill string: U.S. Pat. No. 4,605,268, Russian Federation published patent application 2140527, filed Dec. 18, 1997, Russian Federation published patent application 2040691, filed Feb. 14, 1992, and WO Publication 90/14497A2, Also see: U.S. Pat. No. 5,052,941, U.S. Pat. No. 4,806,928, U.S. Pat. No. 4,901,069, U.S. Pat. No. 5,531,592, U.S. Pat. No. 5,278,550, and U.S. Pat. No. 5,971,072.\nThe U.S. Pat. No. 6,641,434 describes a wired drill pipe joint that was a significant advance in the wired drill pipe art for reliably transmitting measurement data in high-data rates, bidirectionally, between a surface station and locations in the borehole. The '434 Patent discloses a low-loss wired pipe joint in which conductive layers reduce signal energy losses over the length of the drill string by reducing resistive losses and flux losses at each inductive coupler. The wired pipe joint is robust in that it remains operational in the presence of gaps in the conductive layer. The performance attendant these and other advances in the drill string telemetry art provides opportunity for innovation where prior shortcomings of range, speed, and data rate have previously been limiting on system performance.\nWhen a wired drill pipe system is used, it is necessary to have a communication link between the topmost wired drill pipe and a surface processor (which, inter alia, typically performs one or more of the following functions: receiving and/or sending data, logging information, and/or control information to and/or from downhole and surface equipment, performing computations and analyses, and communicating with operators and with remote locations). Various approaches have been suggested, some of which are summarized in U.S. Pat. No. 7,040,415, including use of a slip ring device, and use of rotary electric couplings based on induction or so-called transformer action. A slip ring (also known as brush contact surfaces) is a well known electrical connector designed to carry current or signals from a stationary wire into a rotating device. Typically, it is comprised of a stationary graphite or metal contact (a brush) carried in a non-rotating component which rubs on the outside diameter of a rotating metal ring (e.g., carried on the upper portion of a kelly joint). As the metal ring turns, the electrical current or signal is conducted through the stationary brush to the metal ring making the connection.\nRotary electrical couplings based on induction (transformer action), known, as rotary transformers, provide an alternative to slip rings and contact brushes based upon conduction between rotating and stationary circuitry, so no direct contact is necessary. The transformer windings comprise a stationary coil and a rotating coil, both concentric with the axis of rotation. Either coil can serve as the primary winding with the other serving as the secondary winding.\nThese types of approaches for surface communication have certain limitations and drawbacks attendant the use of complex electromechanical structures, and it is among the objects of the present invention to provide a system for bidirectional communication of signals between the topmost wired drill pipe and a surface processor, with improved efficiency and reliability.\nA further aspect of the drilling and measurement art that is addressed herein relates to safety at the wellsite, and the problem of powering a rotating assembly, at a location that may be classified as a hazardous area, without the use of power carrying wires. Existing techniques have certain limitations. For example, mud turbines, which are powered by the moving drilling fluid, are relatively complex and expensive to build and to maintain. The use of ordinary batteries can be problematic when the drilling operation must be interrupted for battery replacement. It is accordingly among the further objects hereof to provide a safe, efficient, and reliable source of electric power in conjunction with the rotating drill string."}
{"text": "The present disclosure relates to an image processing device, an image forming apparatus, and a computer-readable non-transitory storage medium with an image processing program stored thereon and, more particularly, to image processing in which an amount of toner used in printing is reduced.\nIn image forming apparatuses such as a copier and a printer, technology for performing image processing in which image data serving as a printing target is reduced for the purpose of n-up printing in which printing is performed on one sheet of recording paper on which multiple pages of a document are collected is used. When copies of the multiple pages are printed by the n-up printing, recording paper can be saved. Further, in the image forming apparatuses, technology for performing toner-saving printing by performing image processing in which image data is masked using a preset mask pattern to reduce an amount of toner used when an image of a printing target is printed is also used."}
{"text": "This disclosure generally relates to online digital content streaming, and particularly to an adaptive streaming system that generates enhanced media manifest files for identifying and streaming a live or on-demand media file (e.g., a video file) representations with various video qualities to clients having varying delivery bandwidth and computing processing power.\nAn online system allows its users to connect to and communicate with other online system users. For example, an online system allows a user to encode media files, such as images or videos, and share the media files with other online system users. An online system user may also request to view media files (e.g., stream a video on demand) encoded by the online system user, or by other online system users, encouraging user engagement with the online system. However, this user engagement is decreased with increased platform latency as online system users attempt to stream a media file outside of available bandwidth or view a popular media file.\nTo enable adaptive streaming of multimedia content over HTTP, when an input media file (e.g., a video) is encoded, multiple levels of compression can be applied to the input file producing a plurality of compressed versions (e.g., resolution of 360p, 480p, 720p, 1080p, etc.), or “representations,” of the input file for streaming on a client device. This allows an online system user to select the representation of a media file for streaming that is within the online system user's available bandwidth. However, if the online system user attempts to stream a media file at a representation that requires greater bandwidth than a client device or network connection will allow, the online system user experiences platform latency.\nExisting solutions of adaptive streaming over HTPP divide a media file into segments according to a fixed duration (e.g., 10 second segments) and cache links to these segments (e.g., universal resource locators (URLs)) across computer servers within a given content delivery network (CDN). Each URL containing a segment of the media file is listed on the manifest file of its corresponding media file, in addition to the segment's start points and end points. This allows an online system user to dynamically switch between representations to accommodate available bandwidth. However, creating a manifest file entry for each segment of a media file causes its manifest file to increase in size. Increased manifest file size complicates file administration and decreases manifest file storage efficiency for streaming service, especially live video streaming. In addition, larger manifest files increase platform latency as the manifest file must be downloaded to a client device before the media file may be streamed by an online system user.\nThe latency for fetching the proper media file and/or segments can be compounded by network congestion due to user traffic directed to popular media within a CDN. When a computer server within a CDN has a popular (e.g., trending) media, it is very likely that a large number of user requests for the popular media will be directed to the computer server. This can cause a delay in the delivery of various representations of a requested media file and delay stream traffic.\nFurthermore, for existing solutions, the delay in stream traffic is due, in part, to the initialization process that a streaming system must complete before a client device can commence streaming a media file, such as properly initializing a manifest file associated with a representation of a media file. For example, a manifest file associated with a representation of a media file generally contains information describing an initialization segment (init), which contains data that does not change between segments, a segment index (sidx), which describes byte ranges within a segment, and media segment information. However, using (init+sidx) format in a manifest file to guide a client device on the amount of data to fetch from a storage device is inefficient as it requires the client device to perform a series of fetches in order to locate the first segment of the media file."}
{"text": "The present invention relates generally to immunogenic agents and to agents which enhance the immune response to a selected antigen. In particular, the invention pertains to the use of protein particles as antigens to elicit cellular immune responses.\nNumerous vaccine formulations which include attenuated pathogens or subunit protein antigens, have been developed. Conventional vaccine compositions often include immunological adjuvants to enhance cell-mediated and humoral immune responses. For example, depot adjuvants are frequently used which adsorb and/or precipitate administered antigens and which can retain the antigen at the injection site. Typical depot adjuvants include aluminum compounds and water-in-oil emulsions. However, depot adjuvants, although increasing antigenicity, often provoke severe persistent local reactions, such as granulomas, abscesses and scarring, when injected subcutaneously or intramuscularly. Other adjuvants, such as lipopolysacharrides, can elicit pyrogenic responses upon injection and/or Reiter\"\"s symptoms (influenza-like symptoms, generalized joint discomfort and sometimes anterior uveitis, arthritis and urethritis). Saponins, such as Quillaja saponaria, have also been used as immunological adjuvants in vaccine compositions against a variety of diseases.\nMore particularly, Complete Freund\"\"s adjuvant (CFA) is a powerful immunostimulatory agent that has been successfully used with many antigens on an experimental basis. CFA includes three components: a mineral oil, an emulsifying agent, and killed mycobacteria, such as Mycobacterium tuberculosis. Although effective as an adjuvant, CFA causes severe side effects primarily due to the presence of the mycobacterial component, including pain, abscess formation and fever. CFA, therefore, is not used in human and veterinary vaccines.\nIncomplete Freund\"\"s adjuvant (IFA) is similar to CFA but does not include the bacterial component. IFA, while not approved for use in the United States, has been used elsewhere in human vaccines for influenza and polio and in veterinary vaccines for rabies, canine distemper and foot-and-mouth disease. However, evidence indicates that both the oil and emulsifier used in IFA can cause tumors in mice.\nDespite the presence of such adjuvants, conventional vaccines often fail to provide adequate protection against the targeted pathogen. In this regard, there is growing evidence that vaccination against intracellular pathogens, such as a number of viruses, should target both the cellular and humoral arms of the immune system. More particularly, cytotoxic T-lymphocytes (CTLs) play an important role in cell-mediated immune defense against intracellular pathogens such as viruses and tumor-specific antigens produced by malignant cells. CTLs mediate cytotoxicity of virally infected cells by recognizing viral determinants in conjunction with class I MHC molecules displayed by the infected cells. Cytoplasmic expression of proteins is a prerequisite for class I MHC processing and presentation of antigenic peptides to CTLs. However, immunization with killed or attenuated viruses often fails to produce the CTLs necessary to curb intracellular infection. Furthermore, conventional vaccination techniques against viruses displaying marked genetic heterogeneity and/or rapid mutation rates that facilitate selection of immune escape variants, such as HIV or influenza, are problematic. Accordingly, alternative techniques for vaccination have been developed.\nParticulate carriers with adsorbed or entrapped antigens have been used in an attempt to circumvent these problems and in attempts to elicit adequate immune responses. Such carriers present multiple copies of a selected antigen to the immune system and promote trapping and retention of antigens in local lymph nodes. The particles can be phagocytosed by macrophages and can enhance antigen presentation through cytokine release. Examples of particulate carriers include those derived from polymethyl methacrylate polymers, as well as polymer particles derived from poly(lactides) and poly(lactide-co-glycolides), known as PLG. While offering significant advantages over other more toxic systems, antigen-containing PLG particles suffer from some drawbacks. For example, large scale production and manufacturing of particulate carriers may be problematic due to the high cost of the polymers used in the manufacture the particulate carriers.\nLiposomes have also been employed in an effort to overcome these problems. Liposomes are microscopic vesicles formed from lipid constituents such as phospholipids which are used to entrap pharmaceutical agents. Although the use of liposomes as a drug delivery system alleviates some of the problems described above, liposomes exhibit poor stability during storage and use, and large scale production and manufacturing of liposomes is problematic.\nInternational Publication No. WO 98/50071 describes the use of viral-like particles (VLPs) as adjuvants to enhance immune responses of antigens administered with the VLPs. St. Clair et al. describe the use of protein crystals to enhance humoral and cellular responses. (St. Clair, N. et al, Applied Biol. Sci., 96:9469-9474, 1999).\nDespite the above described adjuvant and antigen-presentation systems, there is a continued need for effective, safe and cost-efficient vaccines with improved purity, stability and immunogenicity.\nThe inventors herein have found, surprisingly, that protein particles are self-sustaining immunogenic agents which produce cellular immune responses. In particular, the active ingredient is also the delivery system, i.e., the protein particles serve as the antigen and the delivery system. Additionally, the inventors have discovered that the protein particles have several advantages (i) the ease of manufacture, (ii) they are more cost-effective to manufacture than existing agents, (iii) they provide for superior immune responses, and (iv) they have reduced toxicity and eliminate the undesirable side-effects observed with other vaccine formulations. Accordingly, then, the invention is primarily directed to the use of such protein particles as antigens.\nIn one embodiment, the invention is directed to an immunogenic composition comprising selected first antigen and a pharmaceutically acceptable excipient, wherein the selected first antigen is a protein particle, and further wherein the protein particle antigen is capable of producing a cellular immune response. In preferred embodiments, the protein particle is formed from a protein selected from the group consisting of a viral, a fungal, a bacteria, an avian or a mammalian protein. In more preferred embodiments, the protein is herpes simplex virus type 2 glycoprotein B (HSV gB2), hepatitis C virus (HCV) or a human immunodeficiency virus (HIV) protein.\nIn another embodiment, the immunogenic composition further comprising an adjuvant, wherein the adjuvant is encapsulated within, adsorbed or conjugated on to, or mixed with the protein particle.\nIn an additional embodiment, the immunogenic composition further comprises a second antigen, wherein the second antigen is distinct from the first antigen, i.e. the protein particle. The second antigen may be a soluble or neutralizing antigen, it may be conjugated on to the protein particle, or it may be associated with a carrier (for example, the second antigen may be encapsulated within, adsorbed or conjugated on to, or mixed with the carrier). In certain preferred embodiments, the carriers include, but are not limited to proteins, polysaccharides, polylactic acids, polyglycollic acids, polymeric amino acids, amino acid copolymers, lipid aggregates (such as oil droplets or liposomes), polymeric particulate carriers, and inactive virus particles. In more preferred embodiments, the carriers comprise a poylmeric particle, wherein the polymeric particle comprises a polymer selected from the group consisting of a poly(xcex1-hydroxy acid), a polyhydroxy butyric acid, a polycaprolactone, a polyorthoester, and a polyanhydride.\nIn an alternative embodiment, the invention is directed to an immunogenic composition comprising a selected first antigen and a pharmaceutically acceptable excipient, wherein the selected first antigen is a protein particle, and further wherein the protein particle is produced by a process comprising the steps of:\n(a) providing an aqueous solution of a protein;\n(b) adding a precipitation agent to the aqueous solution of the protein and stirring the resulting mixture to form the protein particle;\n(c) stabilizing said protein particle by a stabilizing treatment; and\n(d) recovering the protein particles from the aqueous solution.\nIn an alternative embodiment, the aqueous solution is step (a) further comprises an acid, wherein the acid is acetic acid, glycolic acid, hydroxybutyric acid, hydrochloric acid or lactic acid. In preferred embodiments, the precipitation agent comprises oils, hydrocarbons or coacervation agents. In additional preferred embodiments, the stabilizing treatment comprises heat treatment or by treatment with a chemical cross-linking agent.\nIn preferred embodiments, the protein particle is capable of producing a cellular immune response; and is formed from a protein selected from the group consisting of a viral, a fungal, a bacterial, an avian or a mammalian protein. In more preferred embodiments, the protein is herpes simplex virus type 2 glycoprotein B (HSV gB2), hepatitis C virus (HCV) or a human immunodeficiency virus (HIV) protein. In certain preferred embodiments, the cellular immune response can be a cytotoxic-T lymphocyte (CTL) response. In another embodiment, the immunogenic or vaccine composition further comprises an adjuvant and/or a second antigen as described above, wherein the protein particle is capable of functioning as an antigen and/or an adjuvant.\nIn another embodiment, the subject invention is directed to a method for producing a cytotoxic-T lymphocyte (CTL) response in a vertebrate subject comprising administering to the vertebrate subject an immunogenic or vaccine composition comprising the protein particle as described above. The protein particle is administered in an amount effective for eliciting a cytotoxic-T lymphocyte (CTL) response in the vertebrate subject. The protein particles can be co-administered to the subject prior or subsequent to, or concurrent with, an adjuvant and/or a second antigen.\nIn another embodiment, the invention is directed to a method of immunization which comprises administering to a vertebrate subject a therapeutically effective amount of the immunogenic or vaccine composition comprising the protein particle as discussed above.\nThese and other embodiments of the present invention will readily occur to those of ordinary skill in the art in view of the disclosure herein."}
{"text": "A cooling arrangement for a CPU is known from U.S. Pat. No. 6,549,408 B2, BERCHOWITZ. Disposed on this arrangement is an evaporator, and coolant in this evaporator is evaporated by the heat of the CPU and thereby flows upward in a riser conduit to a condenser that is arranged at a higher geodetic level. The vapor is cooled and liquefied there, and flows in response to gravity through a downpipe back to the evaporator, where it is evaporated again.\nThe advantage of such an arrangement is that it operates quietly; a disadvantage, however, is that, according to the second law of thermodynamics, it is impossible in this case to achieve a coolant temperature lower than the ambient temperature, since heat flows only from a higher-temperature medium to a lower-temperature medium. What is obtained instead is a coolant temperature that, because of the temperature gradient necessary for heat transfer, is at least 7 to 10° K (=Kelvin) above ambient temperature; this limits heat dissipation from the CPU, especially on hot days. The overall efficiency of such a system is moreover not good."}
{"text": "LED based luminaires are rapidly replacing both incandescent and fluorescent luminaires for both general lighting and backlighting applications. In large liquid crystal based monitors, and in large solid state lighting applications, such as street lighting and signage, typically the LEDs are supplied in one or more strings of serially connected LEDs, which thus share a common current. A plurality of parallel strings may also be supplied.\nThe power supply which is to drive the LEDs preferably also supplies power to the operating circuitry of the device, thus reducing cost. Typically a single power supply comprising a power transformer with a plurality of secondary windings is utilized, the primary stage of which is controlled by a feedback circuit to provide a fixed direct current (DC) voltage for the operating circuitry of the device through a particular one of the secondary windings.\nA resonant converter is a switching converter that comprises a tank circuit actively participating in determining input to output power flow. Power flow in a resonant converter is typically controlled either by changing the switching frequency, or the duty cycle, or both. In one embodiment two reactive elements, i.e. a capacitor and an inductor form the tank circuit, and such a resonant inverter is known as an LC inverter. A resonant converter comprises a resonant inverter, which has a switching network and a resonant tank circuit, and a rectifier circuit.\nA resonant converter having two inductive elements, and a single capacitor in the tank circuit, where one of the inductive elements is arranged in parallel with the load, and another inductive element is in series with the load and the capacitor, is known as an LLC converter. Advantageously, an LLC converter exhibits a pair of resonant peaks, each associated with a particular one of the inductors. When properly designed, an LLC converter can be simply controlled by adjusting the frequency responsive to an output feedback, as long as the operating frequency is kept between the two resonant peak frequencies. Typically, a drop in output is compensated for by decreasing the operating frequency, and an increase in output is compensated for by increasing the operating frequency.\nIn a typical embodiment, the two inductors are implemented in an integrated transformer having a leakage inductance where the inductance of the primary winding acts as the parallel inductive element and the leakage inductance acts as the series inductive element. The transformer further enables scaling of the design output voltage based on the turns ratio of the primary and secondary windings of the integrated transformer.\nIn order to reduce cost, it is desired to have a single converter provide drive for both the LEDs and for the operating circuits of the device. Since the voltage for the operating circuits of the device must be well regulated, the LED drive voltage is not well regulated. One solution of a circuit 10 for driving at least one LED luminaire offered by the prior art is illustrated in FIG. 1. Circuit 10 comprises: a power source 20; a resonant mode controller 30, optionally comprising an LLC controller; a converter 40 comprising a bridge circuit 50, a primary side capacitance element CP and a transformer 60; a pair of unidirectional electronic valves D1; a unidirectional electronic valve D2; an inductance element 70; a capacitance element C1; a capacitance element C2; an electronically controlled switch SS; a plurality of LED luminaires, denoted respectively L1, L2 and L3; a plurality of electronically controlled switches SL; a plurality of sense resistive elements RS; an LED controller 80; a pair of unidirectional electronic valves D3; a voltage divider 110; and a reference voltage source 120. Bridge circuit 50 comprises a pair of electronically controlled switches, denoted respectively SB1 and SB2. Transformer 60 comprises: a primary winding 130; and a plurality of secondary windings, denoted respectively 140, 150 and 160.\nIn one embodiment, primary side capacitance element CP, capacitance element C1 and capacitance element C2 are each implemented as a capacitor, and are described herein as such. In another embodiment, unidirectional electronic valves D1, D2 and D3 are each implemented as a diode, and are described herein as such. In one embodiment, inductance element 70 is implemented as an inductor, and is described herein as such. In another embodiment, each sense resistive element RS is implemented as a resistor, and is described herein as such. In one embodiment, electronically controlled switches SB1, SB2, SS and SL are each implemented as an n-channel metal-oxide-semiconductor field-effect-transistor (NFET), and are described herein as such.\nThe output of power source 20 is coupled to the drain of NFET SB1 and the return of power source 20 is coupled to the source of NFET SB2. The gates of NFETs SB1, SB2 are each coupled to a respective output of resonant mode controller 30. The source of NFET SB1 is coupled to the drain of NFET SB2 and a first end of primary side capacitor CP. A second end of primary side capacitor CP is coupled to a first end of primary winding 130 of transformer 60 of converter 40. A second end of primary winding 130 is coupled to the return of power source 20. A first and second end of secondary winding 140 of transformer 60 are each coupled to the anode of a respective one of pair of diodes D1 and the center tap of secondary winding 140 is coupled to a common potential.\nThe cathodes of diodes D1 are each coupled to a first end of inductor 70 and to a first end of capacitor C1. The second end of capacitor C1 is coupled to the common potential. The second end of inductor 70 is coupled to the drain of NFET SS and the anode of diode D2. The source of NFET SS is coupled to the common potential and the gate of NFET SS is coupled to an output of LED controller 80. The cathode of diode D2 is coupled to a first end of capacitor C2 and the anode end of each of LED luminaires L1, L2, L3. The second end of capacitor C2 is coupled to the common potential. The cathode end of each of LED luminaires is coupled to the drain of a respective NFET SL and to a respective input of LED controller 80. The source of each NFET SL is coupled to a first end of a respective sense resistor RS and a respective input of LED controller 80. The second end of each sense resistor RS is coupled to the common potential and the gate of each NFET SL is coupled to a respective output of LED controller 80.\nSecondary winding 150 is coupled to a respective load (not shown). Each end of secondary winding 160 is coupled to the anode of a respective diode D3 and the center tap of second winding 160 is coupled to the common potential. The cathode of each diode D3 is coupled to a load (not shown) and a first end of voltage divider 110. A second end of voltage divider 110 is coupled to the common potential and a division junction of voltage divider 110 is coupled to a respective input of resonant mode controller 30. The output of reference voltage source 120 is coupled to a respective input of resonant mode controller 30 and the return of reference voltage source 120 is coupled to the common potential.\nIn operation, resonant mode controller 30 is arranged to alternately open and close NFETs SB1 and SB2 such that primary winding 130 is charged when NFET SB1 is closed and discharged when NFET SB2 is closed. Resonant mode controller 30 typically operates at a fixed duty cycle of near 50%, with a variable frequency, as will be described further. The power supplied to transformer 60 is controlled via secondary winding 160 and voltage divider 110. Particularly, when NFET SB1 is closed and primary winding 130 is charging, power is output from secondary winding 160 via a first diode D3. When NFET SB2 is closed and primary winding 130 is discharging, power is output from secondary winding 160 via the second diode D3. The rectified voltage at the cathodes of diodes D3 is supplied to the load and is additionally divided by voltage divider 110. The divided voltage is compared to the reference voltage output by reference voltage source 120. In the event that the divided voltage is higher than the output of reference voltage source 120, resonant mode controller 30 is arranged to increase the switching frequency of bridge circuit 50 thereby reducing the amount of power supplied to secondary winding 160. In the event that the divided voltage is lower than the output of voltage source 120, resonant mode controller 30 is arranged to reduce the switching frequency of bridge circuit 50 thereby increasing the amount of power supplied to secondary winding 160.\nSecondary winding 140 is similarly influenced by the control of resonant mode controller 30. Since the voltage across secondary winding 140 is not independently controlled, the voltage appearing across capacitor C2 needs to be controlled so as to provide an appropriate operating voltage for LED luminaires L1, L2, L3. The operation of inductor 70, NFET SS and diode D2 act as a boost converter to increase the output voltage of secondary winding 140, stored across capacitor C1. Particularly, when NFET SS is closed, inductor 70 is charged by secondary winding 140. When NFET SS is open, capacitor C2 is charged and LED luminaires L1, L2, L3 are powered by the combination of the power supplied by secondary winding 140 and the power stored on inductor 70. LED luminaires L1, L2, L3 are thus powered at a voltage greater than the voltage provided by secondary winding 140. LED controller 80 is arranged to detect the voltage at the cathode end of each LED luminaire L1, L2, L3 and compare the detected voltages to a predetermined reference voltage. In the event that one or more of the detected voltages are lower than the predetermined reference voltage, i.e. the voltage across capacitor C2 is less than the optimal operating voltages of at least one of LED luminaires L1, L2, L3, LED controller 80 is arranged to increase the duty cycle of the boost converter, i.e. increase the percentage of time that NFET SS is closed. The voltage across capacitor C2 thus increases.\nThe current flowing through each of LED luminaires L1, L2, L3 generates a voltage across the respective sense resistor RS, which is detected by LED controller 80. In one embodiment, LED controller 80 is arranged to adjust the pulse width modulation (PWM) duty cycle of each NFET SL to control the current flowing through each LED luminaire L1, L2, L3. In another embodiment, LED controller 80 is arranged to adjust the gate voltage of each NFET SL to thereby adjust the current flowing through the respective one of LED luminaires L1, L2, L3, by increasing the effective voltage drop across the respective NFET SL. Any excess power is dissipated across the NFET SL. LED controller 80 may be a single unit controlling both NFET SS and the respective NFET SLs, or may be separate control units without exceeding the scope.\nAnother solution of a circuit 200 for driving at least one LED luminaire offered by the prior art is illustrated in FIG. 2. Circuit 200 comprises: power source 20; resonant mode controller 30; converter 40 comprising bridge circuit 50, primary side capacitor CP and transformer 60; capacitor C1; diode D2; inductor 70; capacitor C2; LED luminaire L1; NFET SL; sense resistor RS; LED controller 80; voltage divider 110; reference voltage source 120; and a pair of rectifier bridges, denoted respectively 210 and 220. A single LED luminaire is illustrated, however this is not meant to be limiting in any way and any number of LED luminaires may be provided. Bridge circuit 50 comprises NFETs SB1 and SB2. Transformer 60 comprises: primary winding 130; and plurality of secondary windings 140, 150 and 160. In one embodiment, rectifier bridges 210 and 220 are each implemented as a diode bridge, and are described herein as such.\nThe output of power source 20 is coupled to the drain of NFET SB1 and the return of power source 20 is coupled to the source of NFET SB2. The gates of NFETs SB1, SB2 are each coupled to a respective output of resonant mode controller 30. The source of NFET SB1 is coupled to the drain of NFET SB2 and a first end of primary side capacitor CP. A second end of primary side capacitor CP is coupled to a first end of primary winding 130 of transformer 60 of converter 40. A second end of primary winding 130 is coupled to the return of power source 20.\nEach end of secondary winding 140 of transformer 60 is coupled to a respective input of diode bridge 210 and the return of diode bridge 210 is coupled to a common potential. The output of diode bridge 210 is coupled to a first end of capacitor C1, the cathode of diode D2 and a first end of inductor 70. The second end of capacitor C1 is coupled to the common potential. The second end of inductor 70 is coupled to a first end of capacitor C2 and the anode end of LED luminaire L1. The second end of capacitor C2 is coupled to the anode of diode D2, the cathode end of LED luminaire L1 and the drain of NFET SL. The source of NFET SL is coupled to a first end of sense resistor RS and an input of LED controller 80. The second end of sense resistor RS is coupled to the common potential and the gate of NFET SL is coupled to an output of LED controller 80.\nSecondary winding 150 is coupled to a respective load (not shown), or alternatively not provided. Each end of second winding 160 is coupled to a respective input of diode bridge 220 and the return of diode bridge 220 is coupled to the common potential. The output of diode bridge 220 is coupled to a load (not shown) and a first end of voltage divider 110. A second end of voltage divider 110 is coupled to the common potential and a division junction of voltage divider 110 is coupled to a respective input of resonant mode controller 30. The output of reference voltage source 120 is coupled to a respective input of resonant mode controller 30 and the return of reference voltage source 120 is coupled to the common potential.\nIn operation, resonant mode controller 30 is arranged to alternately open and close NFETs SB1 and SB2, typically at a predetermined duty cycle near 50%. Primary winding 130 is charged when NFET SB1 is closed and discharged when NFET SB2 is closed. The voltage output across diode bridge 200 is controlled by resonant mode controller 30, as described above. Particularly, the voltage across secondary winding 160, rectified by diode bridge 220, is supplied to the load, typically across an output capacitor (not shown) and is additionally divided by voltage divider 110. The divided voltage is compared to the voltage output by reference voltage source 120. In the event that the divided voltage is higher than the output of reference voltage source 120, resonant mode controller 30 is arranged to increase the switching frequency of bridge circuit 50 thereby reducing the amount of power supplied to secondary winding 160. In the event that the divided voltage is lower than the output of voltage source 120, resonant mode controller 30 is arranged to reduce the switching frequency of bridge circuit 50 thereby increasing the amount of power supplied to secondary winding 160.\nThe output by secondary winding 140 is similarly impacted by the operation of resonance mode controller 30, and thus the voltage across capacitor C1 changes responsive to changes in the load of secondary winding 160. The operation of inductor 70, diode D2 and NFET SL act as a buck converter to reduce the output voltage of secondary winding 140, stored across capacitor C1 to the appropriate voltage for LED luminaire L1. Particularly, when NFET SL is closed, power is provided to LED luminaire L1 by secondary winding 140 and inductor 70 is charged by secondary winding 140. When NFET SL is open, LED luminaire L1 is powered by the power stored on inductor 70. LED luminaire L1 is thus powered at a voltage less than the voltage provided by secondary winding 140 responsive to the duty cycle of NFET SL as controlled by LED controller 80.\nThe current flowing through LED luminaire L1 generates a voltage across the respective sense resistor RS, which is detected by LED controller 80. LED controller 80 is arranged to adjust the pulse width modulation (PWM) duty cycle of NFET SL to control the current flowing through LED luminaire L1. Additionally, the PWM adjustment of NFET SL adjusts the duty cycle of the buck converter of inductor 70, diode D2 and NFET SL, thus adjusting the voltage provided to LED luminaire L1 accordingly.\nAdvantageously, a resonant LED luminaire driving circuit, such as the above described LLC converter circuits 10 and 200, automatically provides for zero voltage switching for the LLC switching elements. However, the above converter circuits 10, 200 require an additional inductor 70. Additionally, NFETs SS and SL are operated in hard switching mode. What is desired, and not provided by the prior art, is an integrated converter which provides both a regulated voltage for use by operating circuits, and a regulated LED drive voltage, without requiring the additional inductors, dissipative regulators, or hard switching of the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a magnetic linear scale unit for precisely measuring a relative motion stroke between relatively displaceable two mechanical components. More specifically, the invention relates to a structure for conveniently and efficiently mounting the magnetic linear scale unit onto the relative displaceable components, which mounting structure allows precise adjustment of the magnetic scale position when the scale is installed between relatively displaceable components.\n2. Description of the Background Art\nAs is well known, a magnetic linear scale unit is used for monitoring displacement of one mechanical component relative to the other mechanical component in precise manner. In general, such magnetic linear scale requires a substantially high accuracy in measurement of the stroke of displacement between two mechanical components in an order of micron or ten microns. This requires strict accuracy in mounting the magnetic linear scale unit between the mechanical components. Especially, in mounting of an elongated magnetic scale in a form of a rod or ribbon, for example, substantial precision is required so that the axis of the magnetic scale is arranged precisely parallel to the axis of movement of the mechanical components.\nConventionally, the magnetic scale is rigidly mounted in a scale channel which serves not only as a support for the magnetic scale but also as a magnetic shield. Both ends of the scale channel are so designed as to be rigidly fixed onto one of the relatively displaceable two mechanical components. This means that, since the magnetic scale is mounted rigidly in the scale channel, the accuracy of mount of the scale channel directly reflects on the mounting accuracy of the magnetic scale. On the other hand, the mounting surface, on which an assembly of the scale channel and the magnetic scale is to be mounted is not always precisely flat and accurately parallel to the axis of the movement of the mechanical component or components. When the mounting surface is not precisely parallel to the axis of the movement of the mechanical component or components, additional difficulty has arisen in assuring holding precise parallelism of the magnetic scale relative to the moving axis. It has been a usual technique to compensate the obliquity of the mounting surface to insert a kind of spacer between one end of the scale channel and the corresponding section of the mounting surface. This requires substantially high skill in mounting the assembly of the scale channel and the magnetic scale. Further to say, even for the skilled engineer, it is still a difficult job for mounting the scale assembly with satisfactorily high accuracy.\nIn order to solve this problem in the prior art, an improved mounting structure for the magnetic linear scale unit has been proposed in the U.S. Pat. No. 4,649,648, issued on Mar. 17, 1987, to Kazuo NAGAOKA et al, which has been assigned to the common assignee to the present invention. The invention disclosed in the aforementioned United States Patent was also disclosed in the German Patent First Publication (DE-OS) No. 35 30 776. The prior proposed invention, set forth above, is successful in providing satisfactory adjusting ability for the magnetic linear scale unit for providing accuracy of mounting of the magnetic scale precisely parallel to the moving axis of the movable member."}
{"text": "Information technological devices such as personal computers, handheld devices such as mobile phones and portable media players, and processors embedded in various kinds of equipment are becoming ubiquitous in our everyday environment. At the same time, these devices and their associated software are continually increasing in complexity, leading to an increased need for resolution of complex support issues.\nMost computer support today is provided through telephone support. Telephone support is very labor intensive, and the cost per support call is unavoidably high. As a consequence the purchase price of a personal computer is only a small fraction of the total cost of ownership for most organizations.\nAlso, support incidents not only generate direct costs—they also cause indirect cost, both because of the productive time lost for the user with a computer problem, and because other employees may need to spend their time in helping out. Thus there are very significant economic gains to be made if the process of computer support could be automated to a larger extent within companies. Since many users encounter similar problems, in particular in more homogenous software environments such as those often found within large corporations, such automated solutions could be advantageous even if they handle only a subset of all possible problems.\nFurthermore, telephone support is a less than ideal solution from a working environment perspective. Telephone support is often carried out at large call centers which mostly handle simple repetitive tasks, since many users encounter identical problems. Organizations can thus gain an advantage by providing more rewarding tasks for their employees and letting common repetitive tasks be solved by an automated system.\nSeveral systems have been proposed where both problem diagnosis and problem resolution is automated in computer systems. Such systems have often been based on well-known methods from artificial intelligence, such as expert systems or case-based reasoning. Some examples are U.S. Pat. No. 5,107,500, U.S. Pat. No. 5,944,839, U.S. Pat. No. 5,678,002, U.S. Pat. No. 5,983,364, U.S. Pat. No. 6,742,141 B1, U.S. Pat. No. 6,145,096, and U.S. Pat. No. 6,785,834.\nAnother approach towards more cost-effective computer support is to let the end user take an active role in the diagnosis and resolution of problems, so called self-help. Self-help may for take the form of web-based support databases, where the user is provided with a large number of alternative diagnoses and solutions to a technical problem, and has to choose between these himself.\nSelf-help could also involve communicating with a computer program which attempts to interpret information provided by the user, and also analyzes the state of the user's system, in order to guide the user to a solution of the problem. The user may also automatically be connected to a human support technician in cases where the system fails to resolve the problem at hand. Such systems have previously been described, e.g., in U.S. Pat. No. 6,615,240.\nHowever, it can be very difficult for an ordinary computer user to find relevant technical information in a large database, even with the assistance of a self-help system. A large database of computer support information could today consist of 105 to 106 technical articles (an example is Microsoft Technet, which currently contains at least 150,000 articles). This number is likely to increase significantly in the future as systems become more complex, making the task of navigating a comprehensive support database even more difficult.\nIt is therefore essential to minimize the number of questions asked of the user and the number of alternatives presented in a self-help system. This will serve to increase user acceptance of self-help systems, and increase productivity in the support process.\nOne way of addressing this need is to use information from automated diagnosis tools that determine various aspects of the state of the user's system. This information can be used to guide the search process and to reduce information overload by filtering out alternatives that are irrelevant to the current support context. Such systems have been described, e.g., in U.S. Pat. No. 6,658,598.\nHowever, the amount of support information remaining after filtering may still be very large, making search difficult for the user. This is particularly likely to occur when the support interaction involves limited knowledge of the system on the user's part rather than an actual system fault that could be detected automatically. In these cases, an understanding of typical user behavior is also needed to provide an effective guided search process.\nIt would therefore be beneficial to provide a system and a method for assisting the user in the search process so that the effort in finding a solution to a support problem is minimized."}
{"text": "1. Fieid of the Invention\nThe present invention relates to a free form machining tool having a spherical surface machining section on a lower end for machining a freely curved surface.\n2. Description of the Related Art\nFIG. 1A shows a conventional free form machining tool. FIG. 1B typically shows a freely curved surface to be machined by this machining tool. A conventional free form machining tool 1 is, for example, a ball nose grinder or a ball end mill. The machining tool 1 has a spherical machining surface (such as a grinder or a cutting blade) on the lower end thereof (in this drawing). The machining tool 1 is rotated around an axis z. A freely curved surface 2 is, for example, a pressing mold, an aspherical lens or the like. The free form machining tool 1 is rotated around the axis z at high speed, while the lower end is moved along the freely curved surface 2 so as to work (grind or cut) the freely curved surface 2. A repetition of such a work allows the freely curved surface such as a mold or aspherical lens to be freely machined by the machining tool 1.\nThe aforementioned free form machining tool 1 is rotated around the axis z. Thus, the surface to be machined has a peripheral speed of zero (0) on the axis (radius: 0). Therefore, even if a triaxial (X-Y-Z) NC machining apparatus is used, the position of the axis (radius: 0) is a dead spot. Disadvantageously, this position cannot be excellently machined. Thus, the free form machining tool 1 has been heretofore mounted to a multi-axial NC machining apparatus having four or five axes. For machining a portion (for example, a portion A) where the freely curved surface is perpendicular to the z-axis, a program is made so that the axis z of the machining tool 1 may be appropriately angled with respect to the z-axis. However, such programming is complicated and difficult. Furthermore, the multi-axial NC machining apparatus having four or more axes is expensive. Accordingly, there is a problem that this has little general-purpose properties."}
{"text": "Atherosclerotic plaque is known to build up on the walls of arteries in the human body. Such plaque build-up restricts circulation and often causes cardiovascular problems, especially when the build-up occurs in coronary arteries. Other body passages such as the esophagus, ureter and bile ducts, for example, are subject to blockage by tumorous tissue. Accordingly, it is desirable to remove or otherwise reduce plaque build-up and other tissue obstructions from such body passages.\nKnown catheters use laser energy to remove plaque build-up on artery walls. One such known catheter has a laser source and a catheter body. The catheter body has a proximal end and a distal end, or head, and multiple optical fibers extending between the proximal and distal ends. The laser source is coupled to the optical fibers at the proximal end of the catheter body and is configured to transmit laser energy through the optical fibers.\nTo remove an obstruction from a body passage, such as atherosclerotic plaque in an artery, the catheter is positioned in the artery so that the distal end of the catheter is adjacent to the plaque. The laser source is then energized so that laser energy travels through the optical fibers and photoablates the plaque adjacent the distal end of the catheter. The catheter is then advanced further through the artery to photoablate the next region of plaque build-up.\nWhile known laser catheters are generally acceptable for removing small obstructions, such catheters are limited to opening a path the size of the catheter head on each pass through the body passage. The multiple passes which are required for removing larger areas of obstruction increase the possibility of damaging the passage inner wall. In addition, multiple passes increase the possibility that a piece of the obstruction will break free, enter the blood stream and result in vessel blockage. Other known laser catheters are limited by the relative inflexibility of the catheter distal end which may inflict damage to body passage inner walls as the catheter is advanced.\nAccordingly it would be desirable to provide a laser catheter which can remove substantial portions of an obstruction in a single pass. It would also be desirable to provide a laser catheter having a flexible, adjustable distal end which can substantially conform to the inner dimensions of the body passageway to minimize damage to the inner wall. It would be further desirable to provide a laser catheter which can expand and contract during photoablation to increase the area of obstruction which may be photoablated in a single pass through a body passage."}
{"text": "The zona pellucida is an extracellular matrix that is composed of three glycoproteins, which is a major secretory product of growing oocytes and can also be found on ovulated eggs and early embryos. The fertilized ovum is surrounded by this matrix of mucoproteins, which separates the egg from a layer of follicle cells that provide nourishment to the egg.\nThree components of the mouse zona pellucida have been identified (ZP1, ZP2, and ZP3) and contain a signal sequence that directs the protein into a secretory pathway and is cleaved from the mature protein, and a transmembrane domain. In mice these three genes are expressed in a coordinate, oocyte-specific manner during the growth of oogenesis. See, for example, Epifano, O. et al., Development 121: 1947-1956, 1995.\nIn the mouse model of fertilization, the ZP3 protein first binds the sperm, followed by the binding of ZP2. ZP1 then crosslinks ZP3 and ZP2 dimers. The interaction of these three proteins with the sperm result in the induction of the acrosome reaction (Bork, P. et al., FEBS 300:237-240, 1992). Thus the zona pellucida plays a critical role in fertilization and the post-fertilization block to polyspermy.\nThus genes of the zona pellucida (ZP) family and their encoded proteins are critical in understanding the process of fertilization. The identification of additional members of this family is needed to study this process and to develop methods to modulate it. The present invention provides a novel member of the zona pellucida proteins and related compositions whose uses will be apparent to those skilled in the art from the teachings herein."}
{"text": "Presently, the railroad cars used for containing bulk material are large container cars in which the container compartment is the full size of the car. Special equipment is required to fill the container car with the bulk material and remove the bulk material from the car, generally through doors in the bottom of the car. However, many businesses which desire to ship bulk material are not located near the railroad tracks and do not have easy access to the special equipment needed to transfer the bulk material from trucks to the railroad container car. Also, the business receiving the bulk material may likewise not be near the railroad tracks or have easy access to the special equipment required to transfer the bulk material from the railroad container cars to trucks. There are presently on the market trucks for handling bulk material, such as refuse, having a container which can be easily and quickly slid off of or on to a hoist at the back of the truck. One of such trucks is a Dempster Dumpster with a Dinosaur container made by the Dempster Dumpster systems division of Carrier Corporation of Knoxville, Tenn. Another such truck is made by Accurate division of Western Capital Corporation of Williamstown, N.J. Since these trucks can transfer bulk material in containers which can be easily removed from and replaced on the trucks, it would be desirable to have a railroad flatcar which can easily receive one or more of the containers from such tracks and allow the containers to be transferred back to the trucks. This would allow the delivering business to transfer bulk material in the containers by truck to a railroad siding where the containers could be easily and quickly transferred to the flat car. At the other end of the delivery, the containers could be transferred from the flatcar to a similar truck for delivery to the receiving business."}
{"text": "1. Field of the Invention\nEmbodiments disclosed herein relate generally to wellbore fluids that may include silicate-based additives for stabilizing an unconsolidated formation.\n2. Background Art\nHydrocarbon fluids, such as oil and natural gas, and other desirable formation fluids are obtained from a subterranean geologic formation, i.e., a reservoir, by drilling a well that penetrates the formation zone that contains the desired fluid. Once a wellbore has been drilled, the well must be completed. A well “completion” involves the design, selection, and installation of equipment and materials in or around the wellbore for conveying, pumping, or controlling the production or injection of fluids. After the well has been completed, production of the formation fluids can begin.\nWhen the subterranean formation is “soft” or poorly consolidated, small particulates (typically sand) present in the formation may dislodge and travel along with the produced fluid to the wellbore. Production of sand is highly undesirable as it tends to cause erosion of surface and subterranean equipment, and therefore, it must be removed from the produced fluids before they can be processed. In addition, the migrating sand can plug the flow channels in the formation, thereby necessitating other stimulation techniques, such as acid stimulation, to restore the well's performance.\nVarious types of unconsolidated formations include dune sands, alluvial deposits of sand and gravel, and unconsolidated marine deposits. The challenges in drilling and completing wells in these types of formations are to keep the borehole open and prevent caving, and to avoid reducing the hydraulic conductivity of the near-well formation by introducing irrecoverable mud or smearing clay at the well/aquifer interface during the drilling process. Drilling fluids for unconsolidated formations are typically water-based and generally include clean fresh water, water with clay additives, water with polymeric additives and water with a mixture of clay and polymeric additives.\nOne method of controlling loose sands in unconsolidated formations involves placing a filtration bed of gravel near the wellbore in order to present a physical barrier to the transport of unconsolidated formation fines with the production of hydrocarbons. Typically, such so-called “gravel packing operations” involve the pumping and placement of a quantity of a desired particulate into the unconsolidated formation adjacent the wellbore. Such packs are time consuming and expensive to install.\nAnother method used to control loose sands in unconsolidated formations involves consolidating or stabilizing the unconsolidated subterranean producing zones into hard permeable masses by pre-flushing the formation, applying a hardenable resin composition, applying a spacer fluid, applying an external catalyst to cause the resin to set, and applying an afterflush fluid to remove excess resin from the pore spaces of the zones. Such multiple-component applications, however, often result in uncertainty and create a risk for undesirable results. For example, when an insufficient amount of spacer fluid is used between the application of the hardenable resin and the application of the external catalyst the resin may come into contact with the external catalyst in the wellbore itself rather than in the unconsolidated subterranean producing zone. When resin is contacted with an external catalyst an exothermic reaction occurs that may result in rapid polymerization. The polymerization may damage the formation by plugging the pore channels, may halt pumping when the wellbore is plugged with solid material, or may even result in a downhole explosion as a result of the heat of polymerization. Also, using these conventional processes to treat long intervals of unconsolidated regions is not practical due to the difficulty in determining if the entire interval that has been treated with both the resin and the activation agent.\nThese techniques typically involve the injection of a consolidating fluid, such as a resin-based consolidating fluid, through the wellbore and into the formation surrounding the interval of interest. Resin-based consolidating fluids generally include an organic resin, a curing agent, a catalyst and an oil wetting agent. The resin system hardens in the formation, thereby consolidating it. Examples of such resin-based consolidating fluids and methods for using them are described in, for example, U.S. Pat. Nos. 4,291,766; 4,427,069; 4,669,543; 5,199,492; and 5,806,593. Resin-based consolidation systems may be complicated to apply, especially those involving multiple treatment stages, and the treatment results may be erratic. When the individual components of the consolidating fluid are pumped at different stages into the formation they may or may not come together in the right order, or in the right amounts, or they may not even come together at all. And, even when they do come together, good mixing of the components is not assured, helping to explain the erratic and unreliable results that operators have experienced using such multi-stage consolidating fluids.\nIn an effort to improve performance, other well treatments have been proposed which use inorganic systems, specifically the use of components which form silica gels, to modify the formation and thereby reduce the production of formation fines. For example, U.S. Pat. No. 3,593,796 describes a multi-stage process in which the following components are injected sequentially into the formation: (1) an aqueous solution containing a silicate adapted to wet the fine sand grain particles, (2) an aqueous solution of a silicate-precipitating agent capable of reacting with the silicate in solution (1) so as to form a solidifying material and therein to bind the fine sand grain particles, and (3) a solution containing an oil-wetting agent. This treatment is designed to immobilize the fine particles in the formation and prevent their migration when subjected to subsequent fluid flow. The patent states that aqueous solutions of alkaline earth metal salts (e.g., calcium chloride), acidic iron salts, and certain other metal salts can be used as the silicate-precipitating agent.\nIn another instance, U.S. Pat. No. 3,741,308 describes a method of converting an unconsolidated sand formation into a consolidated, permeable formation by flowing volumes of aqueous calcium hydroxide (or compounds which hydrolyze or react with each other to form calcium hydroxide) through the pores of the unconsolidated formation. The patent states that the calcium hydroxide solution could be formed by adding sodium hydroxide to a solution of calcium chloride. The patent also states that during the practice of the process the sand particles in the formation become coated with calcium silicates of unknown or indefinite composition, and proposes that the coating cements the individual grains together and increases the structural strength of the sand assemblage.\nYet another approach has been described in two companion cases (U.S. Pat. Nos. 5,088,555 and 5,101,901). In U.S. Pat. No. 5,088,555, a sand consolidation method was described involving sequential injections of (a) an aqueous solution of an alkali metal silicate and (b) certain organic solutions of a calcium salt (e.g., calcium chloride hydrate or chelated calcium) through perforations in the casing of a borehole. The components of these two solutions are said to react to form a calcium silicate cement with permeability retention characteristics in the formation interval being treated that prevents sand from being produced during the production of hydrocarbon fluids from the well.\nHowever, the use of a silicate-precipitating agent that is in solution may give rise to short gellation times upon contact silicates. Accordingly, there exists a need for consolidation or stabilization methods which allow for longer and/or controllable gellation times."}
{"text": "A network subscriber can use peer-to-peer (P2P) software for the unauthorized sharing of copyright-protected content (e.g., music) over a P2P network. Typically, Internet Service Providers (ISPs) prevent the unauthorized sharing of copyright-protected material by cutting off the subscriber from the network, essentially assigning the subscriber a bandwidth of zero. Therefore, the network subscriber is not only prevented from transferring copyright-protected material over the P2P network, but the network subscriber is prevented from transferring any material over the network.\nIn addition, ISPs are beginning to operate as Media Providers (MP) (e.g., a Music Service Provider (MSP)) to permit network subscribers to legally share unlimited music files among one another within an enclosed network, or “walled garden.” There are no digital rights management (DRM) restrictions attached to the music files being traded within the MP (e.g., MSP). Therefore, various business rules need to be applied to the content streams (flows) over the network to manage and maintain data transfers within a MP network. For example, MP subscribers are authorized to upload music (i.e., share music) within the MP network. Therefore, content streams of copyright-protected music from a MP subscriber to an ISP subscriber outside of the MP network need to be blocked without disabling the permissible content streams from the same MP subscriber to other MP subscribers within the MP network. The network traffic should be monitored to determine the source of the content, the destination of the content, and the content that is being transferred on a particular network stream. In addition, an appropriate action should be taken on a subscriber's particular network stream to modify (e.g., block the stream) based on the source address, destination address, and content being transferred independent of a subscriber's other network streams."}
{"text": "This invention relates to new and useful improvements in a combination night light and sterilizing holder for toilet articles.\nHolders have heretofore been provided for toilet articles which not only provide storage therefor when not in use but also provide means in the holder for sterilizing the toilet articles. Such previous holders have disadvantages one of which is that the toilet articles are not held in a convenient position for insertion and removal. Another disadvantage is that the holders cannot be readily cleaned on the inside and otherwise readily maintained for efficient usage. Still another disadvantage is that the electric circuitry in the sterilizing means is not substantially concealed to prevent tampering therewith by children."}
{"text": "1. Field of the Invention\nThe present invention relates to a backlight unit suitably used in a liquid crystal display (LCD) device, and more particularly to a backlight unit having an improved structure.\n2. Description of the Related Art\nIn a LCD device, there is provided a LCD panel on the front side of the light-emitting surface of a backlight unit. The LCD panel displays an image by performing optical switching for each of pixels arranged in a matrix. There are two types of backlight units including a central-light-source type in which the light source is provided at the rear side of the light-emitting surface and the light emitted from the light source passes directly through the light-emitting surface, and an edge-light-source type in which the light source is provided at the edge of an optical guide plate provided at the rear side of the light-emitting surface and the light emitted from the light source is guided by the optical guide plate to pass through the light-emitting surface.\nThe edge-light-source backlight unit has the advantage that the light source can be disposed at the side edge of the optical guide plate to provide a small thickness for the backlight unit. Examples of the light source of the edge-light-source backlight unit include a linear cold-cathode fluorescent lamp and a chip-type light emitting diode (LED). In recent years, the chip-type LED is the mainstream of the light source due to smaller dimensions and longer lifetime of the light source. Patent Publication JP-2004-186004A describes an edge-light-source backlight unit using such a chip-type LED as the light source.\nIn the backlight unit using a LED as the light source, in order to emit the light as white light, three color LEDS emitting three colors of red, green and blue are provided without using a single white LED, because the typical white LED has an insufficient emission power. In the backlight unit described in the above patent publication, if the optical path from the color LEDs to the light-emitting surface of the backlight unit has an insufficient length, there arises a problem in that the lights of three colors emitted from the LEDs are not well mixed to reach the light-emitting surface. In this case, the backlight unit cannot provide an even-color backlight at the light-emitting surface, especially at the edge portion thereof.\nIn order to solve the above problem involved with the conventional backlight unit including color LEDs, Patent Publication JP-2006-18175A describes provision of an additional light-mixing optical guide plate provided on the rear side of the main optical guide plate. FIG. 4 is a sectional view showing the configuration of the LCD device described in this patent publication. The LCD device, generally designated by 100, includes a backlight unit 11, a LCD panel 12 provided on the front side of the light-emitting surface 11a of the backlight unit 11, and a shield front 13 configuring a housing and defining the outer edge of the screen of the LCD panel 12.\nColor LEDs 21 are provided at the rear side of a rear cover 38 configuring the rear housing of the backlight unit 11. The color lights emitted from the color LEDs 21 are guided by a first reflector 31, a light-mixing optical guide plate 32, a second reflector 33, and a main optical guide plate 34, thereby passing through an optical sheet 36 configuring the light-emitting surface 11a of the backlight unit to reach the LCD panel 12. In this configuration of the LCD device 100, the color lights emitted from the LEDs 21 are mixed sufficiently by the light-mixing optical guide plate 32, thereby achieving an even color light at the light-emitting surface 11a of the backlight unit.\nFor a color LCD device, it is important to maintain the color reproducibility of the display image for a long period of service. In order to maintain the color reproducibility of the display image during the lifetime of the LCD device, the change of color tones due to aged deterioration should be suppressed on the light-emitting surface of the backlight unit. For this purpose, it is important to suppress the rise of temperature of the LED devices due to the heat generated in the LEDs. If the temperature of the color LEDs rises excessively, the color LEDs will be deteriorated, lowering the luminance of the color LEDs. The reduction in the luminance occurs in an amount depending on the characteristic of the respective LEDs, resulting in a change of the color tone of the backlight.\nIn the LCD device 100 described in JP-2006-18175A, in order to suppress the temperature rise of the color LEDs 21, there is provided a heat sink 23 on the rear surface of a wiring board 22 on which the LEDs 21 are mounted, thereby diffusing the heat generated in the color LEDs 21 via the heat sink 23. In order to further enhance the heat dissipation capability of the LCD device, a structure is employed to enlarge the surface area thereof by providing fins 25 on the heat sink 23, as depicted in FIG. 4.\nHowever, for a LCD device for use in the field of photograph, printing etc., which are required of a higher color reproducibility, the temperature rise should be suppressed as much as possible, and thus the heat dissipation capability should be further improved. In this improvement, a larger surface of the heat sink 23 should be avoided because the larger surface inevitably increases the dimensions of the backlight unit 11. Therefore, it is desirable to improve the heat dissipation capability of the backlight unit without increasing the dimensions of the backlight unit 11."}
{"text": "The invention relates to a method and apparatus for the preparation of multiple gauge metal strip by a shaving operation wherein the strip is drawn through the shaving apparatus over an arced surface.\nIt is necessary to provide a multiple gauge thickness in a metal strip in many applications such as the production of copper strip for the formation of electrical conductors and the like. Conventional procedures such as continuous milling have been employed to produce the desired variations in gauge, however, such processes suffer from the disadvantages of being both time consuming as well as generating an unfavorable form of scrap.\nAn additional procedure which has been known in the art to produce multiple gauge metal strip comprises the reduction to gauge by a rolling operation. However, the rolling operation in the production of multiple gauge strip has been found to suffer from the disadvantages of being restricted to certain shapes, and furthermore tend to involve complex and costly tooling. Additionally, it has been found that the product produced by the aforesaid rolling process fails to meet commercial tolerances as well as being free from other structural defects.\nA further method of producing multiple gauge metal strip is by a drawing process. However, certain complications resulted from this process. In the normal drawing process reduction of the thickness in a section results in an increase in section length so that if the thickness of a given shape varies across its width, the drawing process will result in variable changes in length causing non-uniform metal flow and stresses leading to buckling, twisting, tearing and fracture of the workpiece. Though a wide variety of drawing techniques are known, including the employment of the hydrodynamic principle, none have been suggested or would appear to alleviate the aforenoted deficiencies associated with the drawing of complex multiple gauge configurations.\nThe technique of shaving as a form of metal reduction has been known in the art for some time. However, its application has been generally limited to the finishing of materials of uniform cross section by the reduction of the cross sectional area of the workpiece, with the reduction being conducted along the entire surface thereof. This technique is illustrated in U.S. Pat. No. 3,055,102 to Shaw et al. wherein a shaving tool is applied against a rod or bar which reduces the cross sectional area along the entire surface thereof.\nCertain problems would appear to arise if this shaving technique exemplified by Shaw et al. were to be directly applied to the production of multiple gauge strip products from rectangular stock. Specifically, the application of shaving force against only a portion of the total surface of the strip tends to magnify the problems set forth in the aforenoted patent, particularly, the inability to hold the workpiece properly centered with respect to the tool with the result that the workpiece wanders and a wavy or broken surface may result. A further difficulty which would arise and be magnified by the removal of stock from only a portion of the workpiece surface would be the excessive chattering that takes place as the strip would pass through the tool which would appear on the finished product as a torn or galled surface.\nOne known method of producing multiple gauge strip by draw shaving is illustrated in U.S. Pat. No. 3,992,977 to Winter et al. and assigned to the assignee of the present invention and incorporated herein by reference. The method and apparatus for producing multiple gauge metal strip in the aforesaid U.S. patent while found to be superior to previously known methods still was found to suffer from deficiencies. More particularly, the rake angle at which the tool was disposed was found to be limited to a range of 21/2.degree. to about 25.degree. and was preferably about 5.degree. to 18.degree.. Any increase in rake angle above those critical values would cause the workpiece to ride up the tool resulting in excessive strip material being removed. Thus, as a result of the limit on rake angle, the amount of volume removed per pass was also limited. Furthermore, the manner in which the tool was held in position with respect to the anvil and strip was found to produce chatter marks on the shaved portion of the multiple strip product not unlike those chatter marks encountered in the aforesaid machining operations. The chatter marks are undesirable for functional reasons as well as cosmetic reasons. Chatter marks on the metal strip impedes the bonding of semi-conductor chips to the strip as well as inhibiting the metal plating of the strip which may be necessary in the production of electrical conductors and the like. The shaved gauge variation of the strip was found to be .+-. 0.002\" while variations in the transverse positioning of the cut groove were 0.005\". Finally, the finished multiple gauge strip product was observed to have significant scratching on the backside thereof as a result of being drawn over the flat anvil during the shaving process. Furthermore, the shape of the final product exhibited an increase in cross-wise dish with a corresponding increase in cut depth and/or width. This cross-wise dish can be attributed to the insufficient hold-down forces of the strip against the flat anvil."}
{"text": "A signal transfer time measurement apparatus measures the time from when a signal (radio wave) is transmitted from a base station to a terminal to when a signal (radio wave) transmitted from the terminal is received by the base station. Japanese Laid-Open Patent Publication Nos. 9-170364 and 2003-13644 and Japanese National Phase Laid-Open Patent Publication Nos. 2006-512515 and 2008-515315 each describe an electronic key system including such type of a signal transfer time measurement apparatus. When ID verification is performed between a vehicle and an electronic key through wireless communication, the signal transfer time measurement apparatus measures the time a signal is transferred between the vehicle and the electronic key. By measuring the signal transfer time, ID verification may be prohibited when the electronic key is connected to the vehicle using a relay in an unauthorized manner.\nReferring to FIG. 11, a wideband signal may be used as a transmission signal (impulse) transmitted from a base station to a terminal. FIG. 12 illustrates changes in the signal to noise ratio (SNR) when changing the frequency band of the transmission signal from a narrowband to a wideband. As illustrated in FIG. 12, when widening the frequency band of the transmission signal, the power peak value of the transmission signal received by the base station increases, and the difference increases between the power level of the transmission signal and the power level of a noise wave. This improves the SNR, and increases the likeliness of communication establishment."}
{"text": "Manufacturing a printhead that has relatively high resolution and print-speed raises a number of problems.\nDifficulties in manufacturing pagewidth printheads of any substantial size arise due to the relatively small dimensions of standard silicon wafers that are used in printhead (or printhead module) manufacture. For example, if it is desired to make an 8 inch wide pagewidth printhead, only one such printhead can be laid out on a standard 8-inch wafer, since such wafers are circular in plan. Manufacturing a pagewidth printhead from two or more smaller modules can reduce this limitation to some extent, but raises other problems related to providing a joint between adjacent printhead modules that is precise enough to avoid visible artefacts (which would typically take the form of noticeable lines) when the printhead is used. The problem is exacerbated in relatively high-resolution applications because of the tight tolerances dictated by the small spacing between nozzles.\nThe quality of a joint region between adjacent printhead modules relies on factors including a precision with which the abutting ends of each module can be manufactured, the accuracy with which they can be aligned when assembled into a single printhead, and other more practical factors such as management of ink channels behind the nozzles. It will be appreciated that the difficulties include relative vertical displacement of the printhead modules with respect to each other.\nWhilst some of these issues may be dealt with by careful design and manufacture, the level of precision required renders it relatively expensive to manufacture printheads within the required tolerances. It would be desirable to provide a solution to one or more of the problems associated with precision manufacture and assembly of multiple printhead modules to form a printhead, and especially a pagewidth printhead.\nIn some cases, it is desirable to produce a number of different printhead module types or lengths on a substrate to maximise usage of the substrate's surface area. However, different sizes and types of modules will have different numbers and layouts of print nozzles, potentially including different horizontal and vertical offsets. Where two or more modules are to be joined to form a single printhead, there is also the problem of dealing with different seam shapes between abutting ends of joined modules, which again may incorporate vertical or horizontal offsets between the modules. Printhead controllers are usually dedicated application specific integrated circuits (ASICs) designed for specific use with a single type of printhead module, that is used by itself rather than with other modules. It would be desirable to provide a way in which different lengths and types of printhead modules could be accounted for using a single printer controller.\nPrinter controllers face other difficulties when two or more printhead modules are involved, especially if it is desired to send dot data to each of the printheads directly (rather than via a single printhead connected to the controller). One concern is that data delivered to different length controllers at the same rate will cause the shorter of the modules to be ready for printing before any longer modules. Where there is little difference involved, the issue may not be of importance, but for large length differences, the result is that the bandwidth of a shared memory from which the dot data is supplied to the modules is effectively left idle once one of the modules is full and the remaining module or modules is still being filled. It would be desirable to provide a way of improving memory bandwidth usage in a system comprising a plurality of printhead modules of uneven length.\nIn any printing system that includes multiple nozzles on a printhead or printhead module, there is the possibility of one or more of the nozzles failing in the field, or being inoperative due to manufacturing defect. Given the relatively large size of a typical printhead module, it would be desirable to provide some form of compensation for one or more “dead” nozzles. Where the printhead also outputs fixative on a per-nozzle basis, it is also desirable that the fixative is provided in such a way that dead nozzles are compensated for.\nA printer controller can take the form of an integrated circuit, comprising a processor and one or more peripheral hardware units for implementing specific data manipulation functions. A number of these units and the processor may need access to a common resource such as memory. One way of arbitrating between multiple access requests for a common resource is timeslot arbitration, in which access to the resource is guaranteed to a particular requestor during a predetermined timeslot.\nOne difficulty with this arrangement lies in the fact that not all access requests make the same demands on the resource in terms of timing and latency. For example, a memory read requires that data be fetched from memory, which may take a number of cycles, whereas a memory write can commence immediately. Timeslot arbitration does not take into account these differences, which may result in accesses being performed in a less efficient manner than might otherwise be the case. It would be desirable to provide a timeslot arbitration scheme that improved this efficiency as compared with prior art timeslot arbitration schemes.\nAlso of concern when allocating resources in a timeslot arbitration scheme is the fact that the priority of an access request may not be the same for all units. For example, it would be desirable to provide a timeslot arbitration scheme in which one requestor (typically the memory) is granted special priority such that its requests are dealt with earlier than would be the case in the absence of such priority.\nIn systems that use a memory and cache, a cache miss (in which an attempt to load data or an instruction from a cache fails) results in a memory access followed by a cache update. It is often desirable when updating the cache in this way to update data other than that which was actually missed. A typical example would be a cache miss for a byte resulting in an entire word or line of the cache associated with that byte being updated. However, this can have the effect of tying up bandwidth between the memory (or a memory manager) and the processor where the bandwidth is such that several cycles are required to transfer the entire word or line to the cache. It would be desirable to provide a mechanism for updating a cache that improved cache update speed and/or efficiency.\nMost integrated circuits an externally provided signal as (or to generate) a clock, often provided from a dedicated clock generation circuit. This is often due to the difficulties of providing an onboard clock that can operate at a speed that is predictable. Manufacturing tolerances of such on-board clock generation circuitry can result in clock rates that vary by a factor of two, and operating temperatures can increase this margin by an additional factor of two. In some cases, the particular rate at which the clock operates is not of particular concern. However, where the integrated circuit will be writing to an internal circuit that is sensitive to the time over which a signal is provided, it may be undesirable to have the signal be applied for too long or short a time. For example, flash memory is sensitive to being written too for too long a period. It would be desirable to provide a mechanism for adjusting a rate of an on-chip system clock to take into account the impact of manufacturing variations on clockspeed.\nOne form of attacking a secure chip is to induce (usually by increasing) a clock speed that takes the logic outside its rated operating frequency. One way of doing this is to reduce the temperature of the integrated circuit, which can cause the clock to race. Above a certain frequency, some logic will start malfunctioning. In some cases, the malfunction can be such that information on the chip that would otherwise be secure may become available to an external connection. It would be desirable to protect an integrated circuit from such attacks.\nIn an integrated circuit comprising non-volatile memory, a power failure can result in unintentional behaviour. For example, if an address or data becomes unreliable due to falling voltage supplied to the circuit but there is still sufficient power to cause a write, incorrect data can be written. Even worse, the data (incorrect or not) could be written to the wrong memory. The problem is exacerbated with multi-word writes. It would be desirable to provide a mechanism for reducing or preventing spurious writes when power to an integrated circuit is failing.\nIn an integrated circuit, it is often desirable to reduce unauthorised access to the contents of memory. This is particularly the case where the memory includes a key or some other form of security information that allows the integrated circuit to communicate with another entity (such as another integrated circuit, for example) in a secure manner. It would be particularly advantageous to prevent attacks involving direct probing of memory addresses by physically investigating the chip (as distinct from electronic or logical attacks via manipulation of signals and power supplied to the integrated circuit).\nIt is also desirable to provide an environment where the manufacturer of the integrated circuit (or some other authorised entity) can verify or authorize code to be run on an integrated circuit.\nAnother desideratum would be the ability of two or more entities, such as integrated circuits, to communicate with each other in a secure manner. It would also be desirable to provide a mechanism for secure communication between a first entity and a second entity, where the two entities, whilst capable of some form of secure communication, are not able to establish such communication between themselves.\nIn a system that uses resources (such as a printer, which uses inks) it may be desirable to monitor and update a record related to resource usage. Authenticating ink quality can be a major issue, since the attributes of inks used by a given printhead can be quite specific. Use of incorrect ink can result in anything from misfiring or poor performance to damage or destruction of the printhead. It would therefore be desirable to provide a system that enables authentication of the correct ink being used, as well as providing various support systems secure enabling refilling of ink cartridges.\nIn a system that prevents unauthorized programs from being loaded onto or run on an integrated circuit, it can be laborious to allow developers of software to access the circuits during software development. Enabling access to integrated circuits of a particular type requires authenticating software with a relatively high-level key. Distributing the key for use by developers is inherently unsafe, since a single leak of the key outside the organization could endanger security of all chips that use a related key to authorize programs. Having a small number of people with high-security clearance available to authenticate programs for testing can be inconvenient, particularly in the case where frequent incremental changes in programs during development require testing. It would be desirable to provide a mechanism for allowing access to one or more integrated circuits without risking the security of other integrated circuits in a series of such integrated circuits.\nIn symmetric key security, a message, denoted by M, is plaintext. The process of transforming M into ciphertext C, where the substance of M is hidden, is called encryption. The process of transforming C back into M is called decryption. Referring to the encryption function as E, and the decryption function as D, we have the following identities:E[M]=CD[C]=M\nTherefore the following identity is true:D[E[M]]=M\nA symmetric encryption algorithm is one where:                the encryption function E relies on key K1,        the decryption function D relies on key K2,        K2 can be derived from K1, and        K1 can be derived from K2.        \nIn most symmetric algorithms, K1 equals K2. However, even if K1 does not equal K2, given that one key can be derived from the other, a single key K can suffice for the mathematical definition. Thus:EK[M]=CDK[C]=M\nThe security of these algorithms rests very much in the key K. Knowledge of K allows anyone to encrypt or decrypt. Consequently K must remain a secret for the duration of the value of M. For example, M may be a wartime message “My current position is grid position 123-456”. Once the war is over the value of M is greatly reduced, and if K is made public, the knowledge of the combat unit's position may be of no relevance whatsoever. The security of the particular symmetric algorithm is a function of two things: the strength of the algorithm and the length of the key.\nAn asymmetric encryption algorithm is one where:                the encryption function E relies on key K1,        the decryption function D relies on key K2,        K2 cannot be derived from K1 in a reasonable amount of time, and        K1 cannot be derived from K2 in a reasonable amount of time.        \nThus:EK1[M]=CDK2[C]=M\nThese algorithms are also called public-key because one key K1 can be made public. Thus anyone can encrypt a message (using K1) but only the person with the corresponding decryption key (K2) can decrypt and thus read the message.\nIn most cases, the following identity also holds:EK2[M]=CDK1[C]=M\nThis identity is very important because it implies that anyone with the public key K1 can see M and know that it came from the owner of K2. No-one else could have generated C because to do so would imply knowledge of K2. This gives rise to a different application, unrelated to encryption—digital signatures.\nA number of public key cryptographic algorithms exist. Most are impractical to implement, and many generate a very large C for a given M or require enormous keys. Still others, while secure, are far too slow to be practical for several years. Because of this, many public key systems are hybrid—a public key mechanism is used to transmit a symmetric session key, and then the session key is used for the actual messages.\nAll of the algorithms have a problem in terms of key selection. A random number is simply not secure enough. The two large primes p and q must be chosen carefully—there are certain weak combinations that can be factored more easily (some of the weak keys can be tested for). But nonetheless, key selection is not a simple matter of randomly selecting 1024 bits for example. Consequently the key selection process must also be secure.\nSymmetric and asymmetric schemes both suffer from a difficulty in allowing establishment of multiple relationships between one entity and a two or more others, without the need to provide multiple sets of keys. For example, if a main entity wants to establish secure communications with two or more additional entities, it will need to maintain a different key for each of the additional entities. For practical reasons, it is desirable to avoid generating and storing large numbers of keys. To reduce key numbers, two or more of the entities may use the same key to communicate with the main entity. However, this means that the main entity cannot be sure which of the entities it is communicating with. Similarly, messages from the main entity to one of the entities can be decrypted by any of the other entities with the same key. It would be desirable if a mechanism could be provided to allow secure communication between a main entity and one or more other entities that overcomes at least some of the shortcomings of prior art.\nIn a system where a first entity is capable of secure communication of some form, it may be desirable to establish a relationship with another entity without providing the other entity with any information related the first entity's security features. Typically, the security features might include a key or a cryptographic function. It would be desirable to provide a mechanism for enabling secure communications between a first and second entity when they do not share the requisite secret function, key or other relationship to enable them to establish trust.\nA number of other aspects, features, preferences and embodiments are disclosed in the Detailed Description of the Preferred Embodiment below."}
{"text": "The invention relates to an electronic device having at least one connector for at least one wired sensor, said device having further a sensor break detection unit. Furthermore, the invention relates to a method for detecting a break of at least one wired sensor. Devices and methods according to the invention may be used in the field of automation technology, e.g. for control of manufacturing machinery or data logging.\nIt is known from the art to detect a sensor break by means of under-voltage detection. The input impedance of the electronic device will act as a pull-down resistor bringing an open circuit to 0 V in the case of loss of connection. Sensor break is assumed if the reading collapses to a value of less than 10% of the full range.\nThese known devices have the disadvantage that low measurement values may be erroneous detected as a sensor break. Furthermore, detecting a sensor break takes several measurement cycles, depending on the input impedance and the capacities involved."}
{"text": "There is a focus detection device known in the related art that adopts the split-pupil phase detection method (see patent literature 1). This focus detection device generates a pair of image signals in correspondence to a pair of images formed with a pair of focus detection light fluxes passing through the exit pupil of an optical system. A correlation operation of the known art is executed by shifting the pair of image signals relative to each other, so as to calculate a correlation value representing a degree of coincidence between a pair of image signal strings resulting from the relative shift. Based upon the correlation value, a shift amount at which the highest agree of coincidence is achieved for the patterns expressed in the pair of image signal strings resulting from the relative shift, is detected as a relative image shift amount representing the extent of relative image shift manifested by the pair of subject images. In addition, the focusing condition of the optical system is detected in correspondence to the image shift amount. The focusing condition of the optical system is indicated by the difference between a predetermined focal plane and a detected image plane, i.e., by a defocus amount."}
{"text": "The invention relates to a process for producing two-layer automotive top coats consisting of a base lacquer layer and a clear lacquer finishing layer.\nTwo-layer automotive top coats are produced by wet-on-wet application of a transparent finishing lacquer layer to a colour- and/or effect-giving base lacquer layer that has been pre-dried by exposure to air. The colour- and/or effect-giving base lacquer layer is produced using either non-aqueous or aqueous base lacquers in the lacqering units of the individual motor vehicle manufacturers. Non-aqueous base lacquers contain organic solvents in an order of magnitude of, for example, from 55 to 80 wt. %, while aqueous base lacquers contain only, for example, from 10 to 25 wt. % organic solvents. For reasons of environmental protection, therefore, there is a trend towards replacing base lacquers based on organic solvents by aqueous base lacquers. For example, older lacquering units in which non-aqueous base lacquers are processed are being replaced by new lacquering units constructed for the processing of aqueous base lacquers.\nThe air-exposure conditions for base lacquer layers in lacquering units for the application of base lacquers based on organic solvents are different from those in lacquering units designed specifically for the application of water-borne lacquers. In lacquering units designed for the application of base lacquers based on organic solvents, the air-exposure conditions that prevail are characterised by short air-exposure times of, for example, from 30 to 180 seconds at air temperatures of, for example, from 20 to 30xc2x0 C. Aqueous base lacquers require longer air-exposure times at higher air temperatures in order to obtain a pre-dried base lacquer film on completion of the air-exposure operation. In lacquering units for the application of aqueous base lacquers, therefore, the air-exposure times are, for example, from 5 to 10 minutes and the air temperatures prevailing during the exposure to air are higher at, for example, from 40 to 80xc2x0 C. For example, for the purposes of adequate pre-drying suitable for the subsequent wet-on-wet application of a clear lacquer coating composition, coating layers applied from aqueous base lacquers are exposed to air first for from 1 to 2 minutes at from 20 to 40xc2x0 C. and then for from 3 to 5 minutes using warm air of from 40 to 80xc2x0 C. The short air-exposure times in lacquering units designed for the application of base lacquers based on organic solvents are the result of the only short overall length of the air-exposure region and the belt speed given by the production target of the vehicle body coating process. In the case of base lacquers based on organic solvents, such short air-exposure times are sufficient to obtain a pre-dried base lacquer film on completion of the air-exposure operation, but that is not the case with aqueous base lacquers. Accordingly, it is not possible to process aqueous base lacquers in lacquering units that allow only short air-exposure times.\nThe object of the invention is to provide a process for producing two-layer motor vehicle top coats by the wet-on-wet application of an aqueous base lacquer coating composition and a clear lacquer coating composition while allowing only short air-exposure times of, for example, from 30 to 180 seconds for the base lacquer layer. In particular, the process is to allow aqueous base lacquers to be processed in automated motor vehicle series lacquering units that are designed for the processing of non-aqueous base lacquers and allow only short air-exposure times for the base lacquer layers.\nThe object is achieved by observing the following conditions in the formulation of aqueous base lacquer coating compositions and observing the following process conditions during the exposure to air of the coating layers applied from such aqueous base lacquer coating compositions.\nAccordingly, the invention provides a process for producing a two-layer colour- and/or effect-giving top coat on motor vehicle bodies in an automatic motor vehicle series lacquering unit, wherein a transparent finishing lacquer layer of a clear lacquer coating composition is applied to a base lacquer layer that has been applied from an aqueous colour- and/or effect-giving base lacquer coating composition and exposed to air, and the two coating layers are together baked, which process is characterised in that there is used an aqueous base lacquer coating composition that contains organic solvents according to a high solid value of from 40 to 70%, and the base lacquer layer, after it has been applied and before the transparent finishing lacquer layer is applied, is exposed to air for from 30 to 180 seconds using circulating air at from 25 to 45xc2x0 C. with an air throughput of from 0.10 to 0.70 m/s, based on the area provided with the aqueous base lacquer layer.\nAccording to the invention, aqueous base lacquer coating compositions are used. For example, in addition to water they contain one or more conventional binders, organic solvents and pigments as well as, optionally, fillers, crosslinking agents, and/or additives conventionally employed in lacquers. The aqueous base lacquer coating compositions used in the process according to the invention have, for example, solids contents by weight of from 15 to 50 wt. %; for aqueous effect base lacquers that value is, for example, preferably from 15 to 30 wt. %, and for aqueous single coloured base lacquers it is preferably higher, for example from 20 to 45 wt. %. The solids content by weight of the aqueous base lacquer coating compositions is formed from the sum of the solids contents of the binders, pigments, fillers, crosslinking agents, and non-volatile additives conventionally employed in lacquers. The weight ratio of pigment to binder in the aqueous base lacquer coating composition is, for example, from 0.05:1 to 3:1; for aqueous effect base lacquers it is, for example, preferably from 0.1:1 to 0.6:1, and for aqueous single coloured base lacquers it is preferably higher, for example from 0.1:1 to 2.5:1, in each case based on the weight of the solids. When calculating the ratio of pigment to binder, the sum of the amounts by weight of colour-giving pigments, effect pigments and fillers is related to the sum of the amounts by weight of binder solids, paste resin solids and crosslinking agent solids in the finished aqueous base lacquer.\nIn addition to the air-exposure conditions for the base lacquer layer which are to be observed according to the invention and which are explained hereinbelow, it is essential to the invention that the aqueous base lacquer coating compositions used in the process according to the invention have a high solid value of from 40 to 70%. In the case of aqueous effect base lacquers, that value is, for example, from 40 to 60%, and in the case of aqueous single coloured base lacquers it is, for example, from 40 to 70%. The high solid value of the aqueous base lacquer coating compositions is calculated according to the formula       solids    ⁢          xe2x80x83        ⁢    content    ⁢          xe2x80x83        ⁢    in    ⁢          xe2x80x83        ⁢          wt      .              xe2x80x83            ⁢      %        ⁢          xe2x80x83        xc3x97    100    ⁢    %                                            solids            ⁢                          xe2x80x83                        ⁢            content            ⁢                          xe2x80x83                        ⁢            in            ⁢                          xe2x80x83                        ⁢                          wt              .                              xe2x80x83                            ⁢              %                                +                                              content          ⁢                      xe2x80x83                    ⁢          of          ⁢                      xe2x80x83                    ⁢          organic          ⁢                      xe2x80x83                    ⁢          solvents          ⁢                      xe2x80x83                    ⁢          in          ⁢                      xe2x80x83                    ⁢                      wt            .                          xe2x80x83                        ⁢            %                              \nThe solvent content and the solids content are therefore to be so matched that a high solid value of from 40 to 70% is obtained according to the above formula.\nThe aqueous base lacquers used in the process according to the invention contain conventional ionically or non-ionically stabilised binder systems. Such systems are preferably anionically and/or non-ionically stabilised. Anionic stabilisation is preferably achieved by means of at least partially neutralised carboxyl groups in the binder, while non-ionic stabilisation is preferably achieved by means of lateral or terminal polyethylene oxide units in the binder. The aqueous base lacquers may be physically drying in nature or crosslinkable with formation of covalent bonds. Aqueous base lacquers that crosslink with formation of covalent bonds may be self-crosslinking systems or systems that crosslink by external means. In the latter case, the aqueous base lacquers may be single- or multi-component.\nThe aqueous base lacquers used in the process according to the invention contain one or more conventional film-forming binders. If the binders are not self-crosslinking or self-drying, they may optionally also contain crosslinking agents. Neither the binder component nor the crosslinking component that may optionally be present is subject to any limitations of any kind. There may be used as film-forming binders, for example, conventional polyester, polyurethane and/or poly(meth)acrylate resins. The choice of crosslinking agents that may optionally be present is not critical; it is dependent, in a manner known to those skilled in the art, on the functionality of the binders.\nThe aqueous base lacquers used in the process according to the invention contain colour- and/or effect-giving pigments and, optionally, fillers. Examples of colour-giving inorganic or organic pigments and fillers are titanium dioxide, micronised titanium dioxide, iron oxide pigments, carbon black, silicon dioxide, barium sulfate, micronised mica, talcum, kaolin, chalk, layered silicates, azo pigments, phthalocyanine pigments, quinacridone pigments, pyrrolopyrrole pigments, perylene pigments. Examples of effect-giving pigments are metal pigments, for example of aluminium, copper or other metals; interference pigments, such as, for example, metal-oxide-coated metal pigments, for example titanium-dioxide-coated aluminium, coated mica, such as, for example, titanium-dioxide-coated mica, graphite effect pigments, plate-like iron oxide, plate-like copper phthalocyanine pigments.\nEffect pigments are generally used in the form of a commercial aqueous or non-aqueous paste; organic solvents and additives, preferably water-dilutable organic solvents and additives, are optionally added thereto, and the whole is then mixed with aqueous binder, with shearing. Powdered effect pigments may first be processed with preferably water-dilutable organic solvents and additives to form a paste.\nColouring pigments and/or fillers may be milled, for example, in a portion of the aqueous binder. Milling may preferably also take place in a special water-dilutable paste resin. Milling may be carried out in conventional units known to those skilled in the art. The remainder of the aqueous binder or of the aqueous paste resin is then added to produce the finished colouring pigment mill base.\nThe aqueous base lacquers used in the process according to the invention may contain further additives conventionally employed in lacquers in amounts conventionally employed in lacquers, for example from 0.1 to 5 wt. %, based on their solids content. Examples of such further additives are neutralising agents, antifoams, wetting agents, adhesion-promoting substances, catalysts, flow agents, anti-pitting agents, light stabilisers and thickeners such as, for example, synthetic polymers having groups that are ionic and/or have an associative action, such as polyvinyl alcohol, poly(meth)acrylamide, poly(meth)acrylic acid, polyvinylpyrrolidone, hydrophobically modified ethoxylated polyurethanes or polyacrylates, crosslinked or uncrosslinked polymer microparticles.\nThe aqueous base lacquers used in the process according to the invention contain organic solvents, the amount of which is such that the aqueous base lacquers have high solid values of from 40 to 70%, for example the aqueous base lacquers contain from 20 to 30 wt. % organic solvents. The composition of the organic solvents in the aqueous base lacquers preferably consists of from 30 to 60 wt. % organic solvents that are low-boiling, for example that boil preferably below 120xc2x0 C., and that are miscible with water readily or without a miscibility gap, for example that have a water-solubility at 20xc2x0 C. of more than 70 g per liter of water. Such organic solvents are preferably chosen from methanol, ethanol, n-propanol, isopropanol, n-butanol, 2-butanol, isobutanol, methoxypropanol, methyl ethyl ketone, acetone or mixtures thereof; n- and iso-propanol are particularly preferred.\nExamples of further organic solvents that may be contained in the aqueous base lacquers used in the process according to the invention and that may account for preferably from 40 to 70 wt. % of the organic solvent composition in the aqueous base lacquer are monohydric alcohols having 5 or more carbon atoms, for example hexanol; ethylene glycol ethers or esters, for example ethylene glycol dimethyl ether, diethylene glycol mono- or di-C1-C6-alkyl ethers, butyl glycol, butyl diglycol, ethyl glycol acetate, butyl glycol acetate; propylene glycol ethers or esters, for example propylene glycol dimethyl ether, dipropylene glycol mono- or di-C1-C6-alkyl ethers, ethoxypropanol, propoxypropanol, butoxypropanol, methoxypropyl acetate, ethoxypropyl acetate; ethylene glycol, propylene glycol, and their dimers or trimers, N-alkylpyrrolidones, for example N-methylpyrrolidone; ketones, for example cyclohexanone; aromatic or aliphatic hydrocarbons, for example toluene, xylene or linear or branched aliphatic C6-C12-hydrocarbons.\nIn the process according to the invention, the aqueous base lacquers are applied to motor vehicle bodies that may consist of one type of substrate or of a plurality of types of substrate joined together in a composite construction. In general, the substrates are of metal or plastics. They are generally pre-coated, that is to say plastics substrates may be provided, for example, with a primer coating of plastics, metal substrates generally have a primer coating applied, for example, by electrophoresis and optionally, in addition, one or more further lacquer layers, such as, for example, a primer surfacer layer. The aqueous base lacquers are applied by spraying in one or more spraying operations in a dry layer thickness of from 10 to 50 xcexcm; in the case of aqueous effect base lacquers, the dry layer thickness is, for example, preferably from 10 to 25 xcexcm, and in the case of aqueous single coloured base lacquers it is preferably higher, for example from 15 to 40 xcexcm.\nThe application of the aqueous base lacquer and the clear lacquer is carried out in the process according to the invention by the known wet-on-wet principle, that is to say the base lacquer layer applied from the aqueous base lacquer coating composition is first pre-dried by exposure to air before the clear lacquer coating is applied. Pre-drying by exposure to air according to the invention is important in order that the finished two-layer coating meets the demands made both of the technological properties, such as, for example, adhesion and resistance to the impact of stones, and of the optical properties, such as, for example, colour shade, the development of special effects, and the appearance.\nIn addition to the above-mentioned conditions to be observed in the formulation of the aqueous base lacquer coating compositions used in the process according to the invention, it is essential to the invention that, after application of the base lacquer layers produced from the aqueous base lacquer coating compositions and before application of the transparent finishing lacquer layer, the base lacquer layers be exposed to air for from 30 to 180 seconds using circulating air at from 25 to 45xc2x0 C. with an air throughput of from 0.10 to 0.70 m/s, based on the area provided with the aqueous base lacquer layer.\nExposure to air of the applied aqueous base lacquer layer takes place in the base lacquer air-exposure zone of the motor vehicle series lacquering and lasts for from 30 to 180 seconds, preferably from 60 to 150 seconds. That time is given, for example, by the overall length of the base lacquer air-exposure zone of the motor vehicle series lacquering unit, which is, for example, from 5 to 15 m, and the belt speed of, for example, from 2 to 6 m/min prevailing therein. Exposure to air takes place under circulating-air conditions at air temperatures of from 25 to 45xc2x0 C., preferably from 30 to 40xc2x0 C. The circulating-air conditions are so chosen that an air throughput, based on the area coated with aqueous base lacquer, of from 0.10 to 0.70 m/s, preferably from 0.15 to 0.60 m/s, is used. The air throughput, based on the area coated with aqueous base lacquer, is calculated as the quotient of the volume of air, in cubic meters, that passes through the air-exposure zone per second and the area, coated with aqueous base lacquer, that is located in the air-exposure zone and is to be exposed to air, for example of the order of from 20 to 150 square meters. The volume of air that passes through the air-exposure zone in the process according to the invention is, for example, from 1 to 2 cubic meters per meter of air-exposure zone and per second. The area, coated with aqueous base lacquer, that is located in the air-exposure zone and is to be exposed to air is calculated from the number of coated motor vehicle bodies present in the base lacquer air-exposure zone at the same time, for example from 1 to 3 vehicle bodies, and the area, in square meters, of the vehicle body in question that is coated with aqueous base lacquer and is to be exposed to air, for example of the order of magnitude of from 15 to 35 square meters in the motorcar sector or from 20 to 65 square meters in the commercial vehicle sector. The expression xe2x80x9carea of an individual vehicle body that is coated with aqueous base lacquer and is to be exposed to airxe2x80x9d means not only the area of an individual vehicle body that is subsequently to be provided with a clear lacquer coating, but also includes any portions of the surface that are not to be covered with clear lacquer, for example in the interior space of the vehicle body.\nIn the process according to the invention, exposure to air is carried out under circulating-air conditions. For example, the procedure is such that the circulating air contains from 5 to 15 g of water per cubic meter. It is also possible for a part, for example from 5 to 20%, preferably from 5 to 10%, of the volume of air passing through the air-exposure zone per second to leave the air-exposure zone as outgoing air and to be replaced by a corresponding amount of fresh air, which is mixed in with the air that is being circulated. The fresh air contains preferably less than 15 g, particularly preferably from 5 to 12 g, of water per cubic meter. The water content of the fresh air that is mixed in may be adjusted by means of conventional dehumidification methods, such as, for example, compression of air and/or condensation or absorption of the water from the air.\nThe circulating air is expediently moved at a flow rate of from 4 to 8 m/s, measured at the object. It is preferably a turbulent flow of air, which is directed from top to bottom and from the sides onto the vehicle body provided with the base lacquer layer to be exposed to air. The air is expediently guided uniformly and perpendicularly onto the base lacquer layer to be exposed to air. The air is expediently drawn off downwards.\nThe air-exposure zone may be operated under constant operating conditions or with a variation of individual or several operating parameters. Variation of operating parameters may be carried out in the form of a continual or sudden change over the entire overall length or over one or more portions of the overall length of the air-exposure zone. The air-exposure zone may be divided into one or more, for example from 1 to 3, zones, which may be separated from one another by locks. However, the variation in the operating parameters during the exposure to air always lies within the limits laid down by the invention for the air-exposure operation seen as a whole. For example, it is possible for the air-exposure zone to be divided into 2 zones, the vehicle bodies coated with aqueous base lacquer being exposed to air first in the first zone at low air temperatures, for example from 25 to 30xc2x0 C., and then in the second zone at higher air temperatures, for example from 30 to 45xc2x0 C.; the air throughput in the two zones may, for example, be chosen to be the same or different, for example lower in the first zone than in the second; for example, the flow rate of the circulating air measured at the object may be above 8 m/s in the second zone. When considering the air-exposure zone or the air-exposure operation as a whole, however, the average air throughput lies within the limits specified by the invention.\nThe aqueous base lacquer layers pre-dried by exposure to air are covered with a conventional chemically crosslinking clear lacquer in a dry layer thickness of, for example, from 30 to 100 xcexcm and baked together therewith at temperatures of, for example, from 100 to 150xc2x0 C. An advantage of the process according to the invention arises from the low air temperatures prevailing during exposure to air of the aqueous base lacquer layers. It is thus possible for the base lacquer layers that have been exposed to air to be covered with a clear lacquer without prior cooling of the vehicle bodies. Suitable clear lacquers are in principle all known clear lacquers or transparent pigmented coating compositions. It is possible to use either solvent-containing single-component (1K) or two-component (2K) clear lacquers, water-dilutable 1K or 2K clear lacquers, powder clear lacquers or aqueous powder clear lacquer slurries. The baking conditions for the two-layer coating consisting of the aqueous base lacquer and the clear lacquer are dependent on the clear lacquer system that is used.\nThe process according to the invention makes it possible to produce two-layer motor vehicle top coats consisting of an aqueous base lacquer layer and a clear lacquer finishing layer by the wet-on-wet process while allowing only short air-exposure times of from 30 to 180 seconds for the base lacquer layer. By means of the process according to the invention it is possible to process aqueous base lacquers in automated motor vehicle series lacquering units that are themselves designed for the processing of non-aqueous base lacquers and allow only short air-exposure times for the base lacquer layers. The complex conversion of such lacquering units that are themselves unsuitable for the processing of aqueous base lacquers can be avoided. Provided that the conditions that are essential to the invention are observed during the formulation of the aqueous base lacquers and the air-exposure conditions for the aqueous base lacquer layers that are essential to the invention are observed, it is even possible to carry out the process according to the invention in such a manner that both the aqueous base lacquers and non-aqueous base lacquers are processed in a motor vehicle series lacquering unit."}
{"text": "1. Technical Field\nThe present disclosure relates to a system and method for planning a pathway through an anatomical luminal network of a patient.\n2. Discussion of Related Art\nDuring a surgical procedure clinicians often use CT images for determining a plan or pathway for navigating through the luminal network of a patient. It is often difficult, however, for the clinician to effectively plan a pathway based on CT images alone, especially in the smaller branches of the bronchial tree where CT images typically do not provide sufficient resolution for accurate navigation.\nTo assist a clinician in planning a pathway through a luminal network, automated pathway planning systems and methods have been implemented that automatically generate a pathway from a target designated on a CT image to the an entry point of a patient, e.g., a patient's mouth, nose, other natural entry points or an artificial entry point such, for example, as an incision. One example of an automated pathway planning system and method can be found in U.S. Pat. No. 8,218,846, the entirety of which is incorporated herein by reference.\nWhen using a fully automated pathway planning system, however, the medical device may reach the end of the pathway in an orientation where the working end of the medical device is not oriented toward the target. In this example, a side of the medical device may be oriented towards the target instead of the working end and it may be difficult or impossible for the clinician to gain access to the target with the working end. In particular, when navigating through the small airways of the bronchial tree it may be difficult or even impossible to flex or turn the working end of the medical device towards a target when the target is located perpendicular to the path of travel of the medical device through the small airways."}
{"text": "As communication technologies evolve, the functionality provided to users may improve. More specifically, Time Division Multiplexing (TDM) networks have historically provided communications functionality to call centers and customers. However, the introduction of Voice over Internet Protocol (VoIP) networks and similar technologies have provided the opportunity for call centers to provide improved services to customers, recording of incoming and outgoing communications, and monitoring of the recorded communications. While the utilization of VoIP networks has expanded the capabilities for call centers, VoIP may utilize protocols and technologies that have not historically been utilized. As such, limitations may exist within a VoIP network.\nAs a nonlimiting example, in many current VoIP environments, recording a communication may include passive taps (passive sniffing), and/or bridge connections. In such configurations, one or more parties may not realize that a recording is taking place. Additionally, determining presence data associated with a user of a communications device in an VoIP environment may be unavailable."}
{"text": "A portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent document or the patent disclosure, as it appears in the Patent and Trademark Office patent file or records, but otherwise reserves all copyright rights whatsoever.\nThe present invention relates to analysis of an image that includes features with a prominent skew or orientation. More specifically, the invention relates to the measurement of an angle of skew of lines in an image, including lines of text.\nPeppers et al., U.S. Pat. No. 4,809,344, describe techniques for preprocessing of character recognition. As shown and described in relation to FIGS. 2, 4, and 5A, masks are used to detect laterally and vertically written documents. Masks could be used that have slits inclined at any angle to discriminate an inclined character string or line of characters. As shown and described in relation to FIGS. 6-8, photosensor arrays have different alignment directions and are operated simultaneously to detect laterally and vertically written documents. Photosensor arrays can also be inclined, as shown and described in relation to FIG. 10, to obtain inclination angle data. Col. 6 lines 39-42 indicate that coarse scanning data may be obtained by the scanner of FIG. 2 such that fine scanning data is obtained by the scanner of FIG. 6.\nThe present invention provides automatic techniques that measure skew direction of lines in an image by obtaining a first approximation, then using data indicating the result of the first approximation to obtain a second approximation more precise than the first. In an image of text, the first approximation can be obtained by finding the direction in which characters are closest together, since intercharacter distances are typically smaller than interline distances. Alternatively, it can be obtained by determining which of a number of directions has the largest variance in black pixels per line, because the lines in the direction of skew either have many black pixels due to characters in a line or few black pixels due to spaces between lines. The second approximation can then be obtained by determining which one of a number of directions close to the first approximation has the largest variance in black pixels per line. Similarly, a third, even more precise approximation can be obtained by determining which of a number of directions close to the second approximation has the largest variance in black pixels per line. More generally, this technique can be extended to obtain a sequence of successive approximations of skew direction, each approximation more precise than the previous one.\nOne aspect of the invention is based on the observation of problems with conventional skew measurement techniques. Some skew measurement techniques are simple but imprecise, providing an angle that is accurate to within several degrees of the actual skew angle. Other skew measurement techniques are more precise, but computationally expensive. In addition, some of the more precise skew techniques are subject to a 90xc2x0 error resulting from text configurations in which a white space extends in a direction perpendicular to the lines of text, such as between columns of text.\nThis aspect is based on the discovery of a skew measurement technique that alleviates these problems by making both a coarse or imprecise skew measurement and a fine or more precise skew measurement. The result of the coarse measurement is used as a starting point for the fine measurement. Because it can start within a few degrees of the actual skew angle based on the coarse measurement, the fine measurement may require less computation and is not subject to the 90xc2x0 error described above. This technique can be repeated on the fine measurement to obtain an even more precise skew measurement. In general, this technique can obtain a sequence of successive approximations of skew direction, each approximation more precise than the previous one.\nThe coarse measurement can be implemented by finding an approximate skew of lines in an image. In an image that includes text, for example, distances between characters can be measured in each of a number of initial directions. If the connected components within the characters include pixels of a first color, such as black, and the background is of a second color, such as white, then these distances can be measured between edge pixels, which can be pixels of either color that have at least one neighbor of the other color. Specifically, distance could be measured starting from a black edge pixel and extending in each direction to the nearest black pixel.\nMeasured distance data indicating distances measured in each direction can then be combined, such as by averaging, to obtain central distance data indicating a central value of distances in each direction, and this data can be used to form a profile indicating central value of distance as a function of direction. The minima separated by 180xc2x0 within the profile indicate the direction of intercharacter spacing within the lines of text, because characters are closer together along a line of text than between linesxe2x80x94in general, the intercharacter distances are smaller than the interline distances. Data indicating the direction of one of these minima can be provided as the first approximation of skew direction.\nThe coarse measurement can alternatively be implemented by measuring variance of pixels in a number of initial directions. For example, for a pivot pixel of the first color, i.e. the color of pixels within connected components, the coarse measurement can obtain the square of the number of first color pixels lying along a line extending through the pivot pixel in each of the initial directions. The results for a number of such pivot pixels can then be combined to obtain a sum of squares for each direction. The initial direction with the largest sum of squares provides a first approximation of the skew direction of the text lines because the variance in the number of black pixels is larger at the skew direction than at other directions. In text, for example, lines in the skew direction will either pass through the center of a line of text, with many black pixels, or through a region between lines of text, with only a few black pixels from ascenders and descenders, resulting in a large variance. Data indicating the initial direction with the largest sum of squares can be provided as the first approximation.\nThe fine measurement can be implemented using the approximate skew direction from the coarse measurement and a number of nearby directions. The nearby directions can cover a range that is a small integer multiple of the angle separating the initial directions. For example, for a pivot pixel of the first color, the fine measurement can obtain the square of the number of first color pixels lying along a line extending through the pivot pixel in each of the nearby directions. The results for a number of such pivot pixels can then be combined to obtain a sum of squares for each nearby direction, and the nearby direction with the largest sum of squares provides a second approximation of the skew direction of the text lines. Data indicating the nearby direction with the largest sum of squares can be provided as the second approximation.\nAn even more precise measurement could be implemented using the second approximation and a number of directions near it.\nA machine implementing the invention can include image input means for receiving an input image and providing input image data and image output means for receiving output image data and providing an output image. A processor can store the input image data in memory, modify it in accordance with the skew direction, and then provide output image data to the output means based on the modified image data. For example, the output image could be a deskewed version of the input image. This technique could be applied in a facsimile machine or another digital copier.\nThe following description, the drawings, and the claims further set forth these and other objects, features, and advantages of the invention."}
{"text": "This invention relates generally to coin operated bulk vending machines of the type commonly encountered for dispensing gun, candy, toys or the like. More particularly, the invention relates to bulk vending machines of the type commonly known as the Northwestern Corporation Model 60, and machines similar thereto marketed by others in the bulk vending industry. The prior vending machines such as the Model 60 generally comprise a base which supports a hollow case which in turn supports a merchandise container resting on the upper part of the case. The case houses a coin operated mechanism which controls the receipt of coins for the dispensed product, and controls the dispensing of the product following receipt of sufficient coinage. A fixed elongated rod extending upwards from the base of the vending machine through the hollow case and thence through the merchandise container and the top cover for the merchandise container, secures the components together. The rod is typically provided with a lock member associated with the top cover that screws onto the upper end of the rod and, when tightened, and the key removed, secures the base, the case and the merchandise container in a tight, secure, vending machine assembly.\nIn such prior vending machines rotation of the rotary coin mechanism handle, positioned at the front of the case, would activate the rotary product dispensing wheel in the merchandise container and deposit the coin that permitted the actuation, in the bottom of the base. In some vending machines, such as the Model 60 above noted, the coins are deposited loosely into the base of the machine and are scooped manually from their position in the base. When the change was to be removed, the locking retainer was unlocked and the combined case and merchandise container were moved upwardly by the operator, exposing the front of the base containing the change. The change was then manually removed from the base by scooping it with the hand. In most situations the base was rigidly secured to a mounting bracket, either as a single machine or oftentimes as one of a large group of machines rigidly fixed together. As a practical matter, it was very difficult to collect the coinage in the manner described above, because a single operator is faced with picking up the combined case and container unit with one hand while scooping coinage from the base member with the other hand. This problem was addressed in U.S. Pat. No. 5,467,858 where a structure was described that provided for pivotal motion of the combined case and merchandise container, when in the uppermost position. There, the base and case unit were constructed to permit such a rotation about, essentially, the vertical center axis of the machine, and then resting the combined case and merchandise container unit on rails that formed a part of the base unit. However, this remedy does not work in situations where a plurality of vending machines are mounted side by side on fixed basis, or in other situations where the machine is positioned in close quarters. In such cases, there simply is not enough room to permit the rotation of the case and container elements relative to the base.\nOther prior art devices provided no separately accessible coin box and required complete disassembly of the case and merchandise container in order to service the vending machine.\nThe present invention provides a bulk vending structure that allows vertical movement of the case and container, combined with a very slight forward movement thereof that does not require more side-to-side space than an unmodified type 60 machine, but does provide access to a separate, easily accessible coin box in the base.\nFurther, as a result of a reconfiguration of a separable coin box positioned in the base, the lifting and resting of the case and merchandise container without turning, provides for simple access to the coin box so that the coin box may be physically removed from the frame without difficulty or, alternatively, the operator can remove coins from the coin box while it is resting in the base. Reconfiguration and rearrangement of the assembly rod through an off center position in the base has provided a coin box of maximum foot print and hence maximum capacity.\nThe resulting vending machine has proven to be substantially more efficient than anything known in the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for diagnosing myopia, and in particular relates to a method for diagnosing myopia and/or myopia related complications by determining SNP (single nucleotide polymorphism) genotype in the BICD1 gene. Further, the present invention relates to a method for screening a material for treating myopia and/or myopia related complications.\n2. Description of the Related Art\nMyopia, also called near- or short-sightedness, is a refractive defect of the eye in which collimated light produces image focus in front of the retina when accommodation is relaxed. Those with myopia see nearby objects clearly but distant objects appear blurred. With myopia, the eyeball is too long, or the cornea is too steep, so images are focused in the vitreous inside the eye rather than on the retina at the back of the eye.\nMyopia is a common eye condition worldwide. The prevalence of the condition varies widely among populations, genders, and ages (Invest Ophthalmol Vis Sci 1997; 38:334-40; Optom Vis Sci 2001; 78:234-9; J Formos Med Assoc 2001; 100:68-91). In the USA, the prevalence of myopia was estimated to be approximately 25% between ages of 12 to 54 years (Arch Ophthalmol 1983; 101:405-7). In the Baltimore Eye Survey, myopia was less common in blacks (19.4%) compared with whites (28.1%) (Invest Ophthalmol Vis Sci 1997; 38:334-40). High myopia (defined as refractive dioptric power≦−5.0 D in this study) accounted for 27% to 33% of all myopic eyes, corresponding to a prevalence of 1.7% to 2% in the general population in the USA (Arch Ophthalmol 1983; 101:405-7). Taiwan is among the highest risk areas in the world for myopia. Using the definition of less than −6.0 D for high myopia, high myopia is much more common in Asia. The percentage of myopia in Taiwan is 18% among Taiwanese school boys and 24% among Taiwanese school girls (J Formos Med Assoc 2001; 100:684-91). The totals are even higher than the 13.1% reported among young men in Singapore (Optom Vis Sci 2001; 78:234-9). Furthermore, the prevalence of myopia is increasing in Taiwan based on two large nationwide surveys (participant number>10,000) conducted in 1995 and 2000.\nHigh myopia is associated with potential blinding conditions such as retinal detachment, macular degeneration, and glaucoma. It has been estimated that 5.6% of blindness among school children in the USA is attributable to myopia. Substantial resources are required for optical correction of myopia such as spectacles, contact lenses, orthokeratology, photorefractive keratectomy and laser in situ keratomileusis (LASIK). However, these corrections do not prevent the ocular complications mentioned above. Furthermore, complications arising from the use of contact lenses (Curr Opin Ophthalmol 1998; 9:66-71), orthokeratology (Cornea 2003; 22:262-4) and surgical procedures (J Refract Surg 2003; 19:S247-9) also impose additional risks to myopes. In the USA, treatment of myopia costs an estimated $250 million per year (Arch Ophthalmol 1994; 112:1526-30).\nWhile studies have found that several risks were attributed to environmental factors, twin studies have indicated a strong genetic influence on myopia with the estimates of heritability ranging from 58 to 90% (Invest Ophthalmol Vis Sci 2001; 42:1232-6; Genet Epidemiol 1988; 5:171-81; Hum Hered 1991; 41:151-6; Br J Ophthalmol 2001; 85:1470-6). Using family data, it has been reported that a family history was a significant risk factor for high myopia (Invest Ophthalmol Vis Sci 2004; 45:3446-52). Several studies also demonstrated a similar finding (Optom Vis Sci 1996; 73:279-82; JAMA 1994; 271:1323-7; Invest Ophthalmol Vis Sci 2002; 43:3633-40; Optom Vis Sci 1999; 76:387-92; Invest Ophthalmol Vis Sci 2004; 45:2873-8). However, while few papers reported identifying susceptible myopia genes, none of the studies have been replicated, thus making the identification highly questionable.\nAccordingly, what is needed in the art, is a method for diagnosing, treating, preventing or ameliorating myopia and/or myopia related complications."}
{"text": "1. Technical Field\nThe present invention relates to a test apparatus and a test method.\n2. Related Art\nA system is known that transmits access requests in packets from a transmitter to a receiver. For example, in a system for testing a device under test, access requests are transmitted in packets from a control apparatus to a test module.\nIn such a system, when packets are jammed in a portion of the transmission path, access requests cannot be transmitted from the control apparatus to the test module. In such a case, the entire system must be reset, and so repair takes a long time. Furthermore, the test module cannot be accessed from the control apparatus side, and so it is difficult to analyze the cause of the jam.    Patent Document 1: International Publication Pamphlet 2008/044421"}
{"text": "There are systems for carrying out transmission/reception of data of audio and video (hereinafter referred to as A/V) (hereinafter referred to as A/V data) between equipments linked by digital network to reproduce such data on the real time basis. The system for manipulating A/V equipment existing at a remote place to receive contents by monitor or speaker near at hand to reproduce it is an example of the systems as described above. As an equipment which manipulates such A/V equipment existing at a remote place (hereinafter referred to as remote output equipment), remote controller having speaker or earphone, etc. is mentioned. It should be noted that there are cases of only audio data, or cases of only video data as data caused to undergo transmission/reception.\nIn the case where packet communication is utilized for communication, the transmitting side encodes A/V data to divide it into frames. Further, the transmitting side supplements headers for exchanging data on network to divided individual frames to transmit them as communication packet. In correspondence therewith, the receiving side removes the headers from the received communication packet to decode the removed remaining data portion to reproduce it. By such a series of processing, at the receiving side, it is possible to reproduce audio and video data on the basis of A/V data reproduced at the transmitting side.\nMeanwhile, in such system, there are instances where information except for A/V data is also caused to undergo transmission/reception together. For example, there is the case where user desires to recognize the state of player by remote output equipment at a remote position while reproducing music by player.\nHere, as the player, there is mentioned a disc recording and/or reproducing apparatus having a function to reproduce recording medium, e.g., so-called CD (Compact Disc), DVD (Digital Versatile Disc) or MD (Mini Disc), etc. In addition, as the state that user desires to know, there are mentioned title of tune that player is reproducing, playing remaining time of tune and/or internal state of player (e.g., playback (reproduction), double speed playback (reproduction), stop, temporary stop, etc.).\nHowever, ordinarily, this function is available only in the system where communication except for A/V data can be carried out in parallel. In this case, mechanism for carrying out transmission/reception of information except for A/V data is independently defined to carry out communication in accordance with the mechanism. Further, it is necessary to mount two communication protocols into the equipment.\nFor this reason, players designed in correspondence with only communication protocol of A/V data in which such a definition is not carried out in advance cannot send, on the real time basis, information such as title of tune being reproduced, playing remaining time of tune and/or internal state of player, etc. to remote output equipment."}
{"text": "A wide variety of devices have been proposed for securing a pair of skis together in running surface or base to base relationship, with the bases spaced apart so as to avoid damaging the running surfaces of the skis. These devices have taken various forms, such for example s the mechanical devices shown in U.S. Pat. Nos. 2,562,178 and 3,626,553, the flexible belt devices shown in U.S. Pat. Nos. 3,947,927 and 4,120,437, and the simple flexible straps of U.S. Pat. Nos. 3,731,348 and 3,841,648.\nThe latter two patents disclose relatively inexpensive ties comprised simply of a length of fabric with appropriately located hook and lop fabric fastening means thereon and arranged so that one end of the tie or strap may be inserted between the bases of the two skis and the remainder of the strap may be wrapped around the two skis to hold them together, the hook and loop fastening means being so located as to secure the other end of the strap to the portion of the strap which underlies said other end when the strap is wrapped about the skis. When not being used to tie the skis together, the flexible straps are simply put in the skier's pocket.\nThe concept is simple and straight forward, but the straps are very difficult to use because the one end of the strap must be manually clamped and held between the skis while the other end is being wrapped around the skis. This requires two hands for use of the strap. Moreover, in cold weather and in the presence of snow and ice, tying of the skis together is difficult to accomplish because the one end of the strap will usually slip out of place before the remainder of the strap can be wrapped around the two skis."}
{"text": "1. Field of the Invention\nThis invention relates to a thermoplastic semiconductor package and to a method of making it, and more particularly concerns attaching a thermoplastic lid or cover to a thermoplastic cavity type semiconductor package body by ultrasonic welding. The semiconductor package may be a Plastic Quad Flat Pack package (PQFP) as shown in FIGS. 1 and 11, Plastic Dual In-Line Package (PDIP) as shown in FIGS. 3 and 13, or a Plastic Pin Grid Array package (PPGA) as shown in FIG. 5, and may include Small Outline Integrated Circuit (SOIC) variations, but is not limited thereto. The semiconductor package may include single or multiple Integrated Circuit (IC) arrangements.\nThe term \"Plastic Cavity Type Semiconductor Package\" refers to any type of semiconductor package that provides for the mounting of an integrated circuit die onto the die pad of the semiconductor package through a recess or cavity in the semiconductor package, subsequently adding electrical interconnections and sealing the electrical components by applying a sealant over the wires and die.\nThe terms \"lid\" or \"cover\" as used in this disclosure are interchangeable, and refer to any component that is attached to the main body of a semiconductor package to protect the internal components of the semiconductor package, such as the integrated circuit, from environmental or physical damage or contamination. The design of the lid may be such to fit within the constraints of a cavity in the semiconductor package body. Alternatively, if desired, the lid may be an integral component of the overall cover of the device. The term lid does not define top or bottom.\nThe term \"ultrasonic welding\" refers to a method of welding elements together by producing an intense, controlled vibration from an ultrasonic welding booster horn and applying it to the elements to be welded. This intense controlled vibration, coupled with an application of a set pressure against the components to be welded, develops frictional heat which causes melting of one or both of the components at the interface between the components, and creates a mechanical bond between the components.\nA plastic cavity type semiconductor package is a practical vehicle for the mounting of a subminiature integrated circuit to printed circuit or wiring board. Plastic cavity-type semiconductor packages contain an integrated circuit which is electrically interconnected to an integral polyester, epoxy or polyimide resin base laminated substrate with an integral electrical circuit, or to an integral conductive lead frame. The substrate circuit or lead frame enlarges the effective pitch, or spacing, between electrical connections of the integrated circuit to a practical, useable dimension found on printed circuit or wiring boards. The method of enlargement may be by a flexible or rigid substrate circuit, a stamped or chemically etched conductive lead frame, or a combination of a substrate circuit and a conductive lead frame.\nThe ways of mounting a plastic cavity type semiconductor package onto a printed circuit or wiring board varies from insertion or through hole mounting to surface mounting.\nThe particular way of mounting dictates the arrangement of the electrical leads that interface with the printed circuit or wiring board.\nThe method of pitch enlargement typically utilized for Plastic Quad Flat Pack (PQFP) and Plastic Dual In-Line Packages (PDIP) ma be either an integral substrate circuit electrically connected to the internal ends of a lead frame, whereby the external ends of the lead frame interface the printed circuit or wiring board, or an integral lead frame whereby the integrated circuit is electrically connected to the internal ends of the lead frame and the external ends interfacing the printed circuit or wiring board. The interface may be a surface mount or insertion type interface.\nThe method of enlargement typically utilized for a Plastic Pin Grid Array, (PPGA) may be an integral substrate circuit electrically connected to the internal ends of a conductive pin, whereby the pins form a matrix pattern and provide electrical connection to the printed circuit or wiring board by inserting the external ends of the conductive pin matrix into the printed circuit or wiring board.\nUpon mounting the integrated circuit or die onto the die pad, a series of wire interconnections, or wire bonds, are made between the integrated circuit and the internal leads of a lead frame, or to a substrate circuit. In applications requiring sealed integrated circuits, a sealant is applied on and about the integrated circuit and the wire bond interconnections.\nPlastic cavity type semiconductor packages are offered in various formats. Plastic Pin Grid Arrays, PPGA, for example, are available in a bare substrate circuit and pin matrix in which all the components compromising the semiconductor package are visible. A cavity, or recess, is designed into the substrate circuit for mounting of an integrated circuit die.\nThe PPGA is also available with the substrate circuit, and the interface of the pin matrix and the substrate circuit, pre-molded or encapsulated with a plastic molding compound. A cavity is designed into this package body for mounting an integrated circuit die and for electrically interconnecting the die to the substrate circuit.\nThe end user provides for the integrated circuit die mounting, the electrical interconnection of the integrated circuit to the substrate circuit and the protection of the integrated circuit by the application of a sealant or lid onto the semiconductor package.\nIn the case of a bare semiconductor package, a typical method of providing protection to the integrated circuit and electrical interconnection against environmental or physical damage involves the total encapsulation of the bare package with a molding compound, or the application of a \"Blob Top\", by placing a suitable quantity of sealant, in the cavity on top of the integrated circuit and electrical interconnections, and curing the sealant by applying heat or ultraviolet rays.\nIn the case of the pre-molded or encapsulated cavity package, a typical method of providing protection to the integrated circuit and electrical interconnection involves applying a \"Blob Top\" or attaching a lid over the cavity. In a typical lidding process, a suitable amount of adhesive is applied about the cavity opening or on the lid or, in some cases, both. The lid and package body are then mated together and allowed to cure by applying heat or ultraviolet rays. The actual application of the adhesive may be by dispensing a liquid material, screen printing a paste material, or applying a preformed adhesive material blank onto the components to be mated.\nThe disadvantages of the current methods of protecting the integrated circuit in a Plastic Cavity-Type Semiconductor Package include the following:\na Adhesive bonding of the lid to the semiconductor package is a time consuming process that involves subjecting the assembly to possible contamination by the adhesive bonding agents. PA1 b. Adhesive bonding agents, and sealing agents, curable by any method, are hydroscopic compounds which absorb moisture, and this moisture is introduced into the cavity and onto the electrical components, causing corrosion and failure. PA1 c. Moisture within a lidded plastic cavity semiconductor package causes an increase in pressure within the cavity when subjected to elevated temperatures, such as soldering. The increase in pressure, due to the expansion of the moisture, raises internal pressures to a level which may cause the lid to separate from the assembly. PA1 d. Encapsulation, or Blob Topping, is a time consuming process which involves subjecting the assembly to possible contamination by the sealing agent, possible failure of the assembly due to moisture introduction, and possible failure of the assembly due to thermal induced stress, and may cause separation or disconnection of the wire bond from the integrated circuit, or separation of the electrical interconnection to the semiconductor package. Heat curing of the sealant also causes thermal induced stress. PA1 e. Over molding semiconductor packages with a plastic mold compound involves subjecting the assembly to possible wire sweeping, which occurs during the injection of the plastic molding compound. Wire sweeping may cause shorting between adjacent wire bonds, or may cause separation or disconnection of the wire bonds from the integrated circuit, or the electrical interconnection to the semiconductor package. PA1 f. Stress critical integrated circuits may not be subjected to multiple cycles of elevated temperature, as the integrated circuit may crack or disconnect from the die mount surface. Failures due to stress cracking are increased by thermal cycling. An example of multiple thermal cycling: Die attaching requires the introduction of heat to cure the bonding agent, and Blob Topping requires heat to cure the sealant, and over molding, which is performed at elevated temperatures. PA1 g. Additional handling of a semiconductor package during the process of adhesive bonding or sealing involves subjecting the assembly to possible contamination or physical damage which may cause failure. PA1 h. A major disadvantage in the assembly of an over-molded or blob topped semiconductor package is noticed during thermal cycling. The differences in the coefficients of thermal expansion, the rate in which a material expands or contracts in response to temperature changes, of the integral components, i.e., the die, the die attach agents, the wire interconnections, the encapsulation material, and so on, are such that when the assembly is subjected to thermal cycling, multiple exposures to hot and cold temperatures, the components expand or contract at different rates. Thermal cycling causes stress and/or strain upon the components due to the differences in those rates and may cause component failure due to wear out, disassociation of surfaces, or fractures of the integral components. Also, the larger dies which are now utilized, include integral fine line geometry. During thermal cycling, the differences in thermal expansion may cause the fine line geometry to crack or separate due to the stresses upon the die surface caused by the expansion of the encapsulation material against the die surface. This is an increasing occurrence due to the increased use of large, fine line geometry, stress critical dies. Another disadvantage of thermal cycling and the problem associated with thermal expansion, is the additional potential for failure due to corrosion or lid separation. As the integral components expand and/or contract during thermal cycling, pockets or voids are created. These pockets or voids allow moisture to gather. Moisture within a semiconductor package causes failures due to corrosion or increase in internal package pressure."}
{"text": "Communication networks increasingly carry information at various data rates. The selected data rate for a given communication can be pre-established between the transmitter and receiver, can be signaled, for example, using handshaking techniques, or can by dynamically determined by the receiver. Data rate detection techniques allow a receiver to determine the rate of incoming data, for example, by examining the received data. Such automatic data rate detection techniques allow a receiver to receive data from a variety of transmitting devices operating at different speeds without having to establish data rates in advance\nSignals arriving at a receiver are typically corrupted by intersymbol interference (ISI), crosstalk, echo, and other noise. In order to compensate for such channel distortions, communication receivers often employ well-known equalization techniques. For example, zero equalization or decision-feedback equalization (DFE) techniques (or both) are often employed. Such equalization techniques are widely-used for removing intersymbol interference and to improve the noise margin. See, for example, R. Gitlin et al., Digital Communication Principles, (Plenum Press, 1992) and E. A. Lee and D. G. Messerschmitt, Digital Communications, (Kluwer Academic Press, 1988), each incorporated by reference herein Generally, zero equalization techniques equalize the pre-cursors of the channel impulse response and decision-feedback equalization equalizes the post cursors of the channel impulse response.\nA communication channel typically exhibits a low pass effect on a transmitted signal. The various frequency components of a signal will thus encounter different attenuation at the output of the channel, with higher frequency components of a transmitted signal being impaired more than lower frequency components. Thus, the impairment of a channel is said to be rate-dependent. As a result, the equalization parameters optimized for one data rate will typically not be applicable for another data rate\nIn the absence of a received signal, the receiver lacks information (data transitions) and cannot sustain a frequency lock. If the equalizer is allowed to train when the signal has been lost, the equalizer will produce invalid updates Likewise, there are a number of predefined patterns that are not sufficiently spectrally rich to provide valid equalization results. For example, many communications systems continuously send an idle pattern to keep the system alive, in a similar manner to a heart beat signal. The idle pattern, however, is not spectrally rich and is therefore not good for equalization.\nA need therefore exists for rate-dependent methods and apparatus for equalizing a channel A further need exists for equalization methods and apparatus that can detect the data rate, and perform equalization based on the detected data rate. Yet another need exists for equalization methods and apparatus that update the equalization parameters only if one or more predefined qualifier conditions, such as a loss of signal, are not present."}
{"text": "Field of the Disclosure\nThis application relates, in general, to a wake system for a watercraft, and more particularly, to a surf wake system for modifying a wake produced by a watercraft travelling through water.\nDescription of the Related Art\nWake surfing has become increasingly popular in recent years because, unlike an ocean wave, a wake produced by a watercraft is on-demand not to mention continuous and endless as long as the watercraft is moving forward. As a watercraft travels through water, the watercraft displaces water and thus generates waves including bow wave and diverging stern waves on both sides of the watercraft. Due to pressure differences, these waves generally converge in the hollow formed behind the traveling watercraft and/or interfere with each other to form a wake behind the watercraft. Such a wake, however, is generally small, choppy or too close to the watercraft to be suitable and safe for water sports, and particularly not suitable for wake boarding or surfing.\nTo facilitate surfing, a wake should be formed away from the stern of the watercraft, for example, about ten feet away, and with a waist-height peak, for example, about three feet or higher. Generally hundreds, and sometimes thousands, of pounds of additional weight or ballast to a rear corner of the watercraft to make the watercraft tilt to one side, displaces more water, and hence generates a larger wake on that side. Such additional weight may be in the form of removable ballast bags, installed ballast tanks or bladders, or passengers positioned to one side of the watercraft, which is primarily used to tip the watercraft to that side. Using such additional weight to produce larger wakes, however, poses several disadvantages. For example, such additional weight may take up significant space and capacity that may otherwise reduce the passenger capacity of the watercraft. Also, such additional weight may unbalance the watercraft creating difficulties in control. Moreover, the additional weight generally must be moved from one side of the water craft to the other in order to generate a wake on the other side of the water craft. Shifting such additional weight may require significant time and effort. For example, filling and emptying ballast tanks to switch from one side to the other may require 20 minutes or more.\nAlternatively, it is known to require extensive modification to a boat hull to promote a proper surf wake. An exemplar of generating a larger wake can be found in a U.S. Pat. No. 6,105,527 to Lochtefeld et al.\nIn light of the foregoing, it would therefore be useful to provide surf wake system that overcomes the above and other disadvantages."}
{"text": "In typical industrial plants and food processing facilities, safety practices such as human hazard reporting and safety observation tours (SOTs) address safety hazards as they are identified and dealt with. However, despite these efforts, accidents, near misses, and contamination issues still occur. Often, safety issues are addressed with quick fixes that result in corrosion of metal, overstressing of electrical systems, and other conditions that eventually lead to equipment malfunctions and new safety hazards. Further, while some facilities may include sensors to identify specific, isolated safety concerns, such as ingress of liquid past a predefined boundary or electrical overloads, complex and slowly developing safety issues at an industrial site may go unnoticed."}
{"text": "The invention relates generally to spherical-shaped semiconductor integrated circuit manufacturing systems, and more particularly, to a fluid exchange system, suitable for use in such manufacturing systems, for purging a first fluid from a process flow comprised of the first fluid and spherical-shaped semiconductors and injecting a second fluid into the process flow.\nConventional integrated circuits, or \"chips,\" are formed from a flat surface semiconductor wafer. The semiconductor wafer is first manufactured in a semiconductor material manufacturing facility and is then provided to a fabrication facility. At the latter facility, several layers are processed onto the semiconductor wafer surface. Once completed, the wafer is then cut into one or more chips and assembled into packages. Although the processed chip includes several layers fabricated thereon, the chip still remains relatively flat.\nA fabrication facility is relatively expensive due to the enormous effort and expense required for creating flat silicon wafers and chips. For example, manufacturing the wafers requires several high-precision steps including creating rod-form polycrystalline semiconductor material; precisely cutting ingots from the semiconductor rods; cleaning and drying the cut ingots; manufacturing a large single crystal from the ingots by melting them in a quartz crucible; grinding, etching, and cleaning the surface of the crystal; cutting, lapping and polishing wafers from the crystal; and heat processing the wafers. Moreover, the wafers produced by the above processes typically have many defects which are largely attributable to the difficulty in making a single, highly pure crystal due to the above cutting, grinding and cleaning processes as well as due to the impurities, including oxygen, associated with containers used in forming the crystals. These defects become more and more prevalent as the integrated circuits formed on these wafers become smaller.\nAnother major problem associated with modern fabrication facilities for flat chips is that they require extensive and expensive equipment. For example, dust-free clean rooms and temperature-controlled manufacturing and storage areas are necessary to prevent the wafers and chips from defecting and warping. Also, these types of fabrication facilities suffer from a relatively inefficient throughput as well as an inefficient use of the silicon. For example, facilities using in-batch manufacturing, where the wafers are processed by lots, must maintain huge inventories to efficiently utilize all the equipment of the facility. Also, because the wafers are round, and the completed chips are rectangular, the peripheral portion of each wafer cannot be used.\nStill another problem associated with modern fabrication facilities is that they do not produce chips that are ready to use. Instead, there are many additional steps that must be completed, including cutting and separating the chip from the wafer; assembling the chip to a lead frame which includes wire bonding, plastic or ceramic molding and cutting and forming the leads, positioning the assembled chip onto a printed circuit board; and mounting the assembled chip to the printed circuit board. The cutting and assembly steps introduce many errors and defects due to the precise requirements of such operations. In addition, the positioning and mounting steps are naturally two-dimensional in character, and therefore do not support curved or three dimensional areas.\nTherefore, due to these and various other problems, only a few companies in the world today can successfully manufacture conventional flat chips. Furthermore, the chips must bear a high price to cover the costs of manufacturing, as well as the return on initial capital and investment.\nIn co-pending U.S. Pat. No. 5,955,776 filed on May 16, 1997, assigned to the same assignee as the present application and hereby incorporated by reference as if reproduced in its entirety, a method and apparatus for manufacturing spherical-shaped semiconductor integrated circuits is disclosed. As disclosed in the aforementioned patent application, the manufacturing process by which a spherical-shaped semiconductor integrated circuit is produced may include a variety of processing steps. Among these are: de-ionized water cleaning, developing and wet etching; diffusion, oxidation and deposition of films; coating; exposure; plasma etching, sputtering and ion implantation; ashing; and epitaxial growth. Many of these steps involve exposing a spherical-shaped semiconductor to a process fluid. For example, a film having a desired composition and thickness may be formed on the surface of the spherical-shaped semiconductors by exposing the spherical-shaped semiconductors to a first process fluid, typically, a gas, having a selected composition at a selected pressure and temperature. To ensure that formation of spherical-shaped semiconductor integrated circuits proceeds as designed, when such a process step is complete, it is often advisable to remove the first process fluid from further contact with the spherical-shaped semiconductors, for example, using a purging process. Failure to do so may result in unwanted chemical reactions, for example, the deposition of an overly thick layer of material, on the surface of the spherical-shaped semiconductors. Additionally, as the manufacturing process for spherical-shaped semiconductor integrated circuits often incorporates a multitude of processing steps such as the ones previously described, to aid in the efficiency of the manufacturing process and/or to assist in the transport of the spherical-shaped semiconductors between processing stations, it is often desirable to replace the removed process fluid with a second process fluid.\nFurther complicating the manufacturing process for spherical-shaped semiconductor integrated circuits are the special handling requirements which must be afforded the circuits during the manufacture thereof. Unlike conventional integrated circuits formed from flat surface semiconductor wafers, the spherical-shaped semiconductor integrated circuits produced from spherical-shaped semiconductors consume the entire surface area of the spherical-shaped semiconductors. Thus, unlike conventional integrated circuits which may be grasped along bottom or side surfaces thereof, grasping or otherwise contacting spherical-shaped semiconductor during the manufacturing process may result in significant damage thereto. Thus, many of the processing techniques used to manufacture conventional integrated circuits are unsuitable for use in the manufacture of spherical-shaped semiconductor integrated circuits.\nThus, it would be desirable to provide a fluid exchange system, capable of purging a first fluid from and injecting a second fluid into a process flow, suitable for use in spherical-shaped semiconductor integrated circuit manufacturing processes. It is, therefore, the object of the invention to provide such a fluid exchange system."}
{"text": "Currently there are various kitchen appliances using electric motors. Some kitchen appliances desire the electric motor to be smaller without reducing its power. One such kitchen appliance is the food processor.\nConventionally, the food processor uses a two-pole universal motor. The stator of the motor has a stator core which includes a yoke with a rectangular configuration and two salient poles which extend inwardly from two opposing sides of the yoke and have stator windings wound thereon. When the stator windings are electrified, two magnetic poles with opposite polarity are formed at the salient poles and two magnetic circuits are formed. Each magnetic circuit passes through two stator poles, the rotor and a side of the yoke so that it has a relatively long path.\nThe present invention aims to provide a new kitchen appliance, especially a food processor which is permitted to be smaller and cheaper."}
{"text": "1. Field of the Invention\nThe present invention pertains to a method for disposing of slurried and sometimes toxic solids wastes in hydraulically fractured subterranean earth formations and in such a way that the waste material is confined in the fracture.\n2. Background\nCertain types of solid waste material may be permanently stored or disposed of by forming a slurry of solids particulates of the waste with a carrier liquid, reducing the particle size of the solids and then injecting the slurry into a subterranean earth formation. U.S. Pat. No. 5,129,469 to J. E. Jackson, issued Jul. 14, 1992, assigned to the assignee of the present invention, and incorporated herein by reference describes a system for disposing of drill cuttings from drilling wells in earth formations by combining the cuttings with a carrier liquid, grinding the cuttings to a predetermined particle size and injecting the viscous slurry of cuttings and carrier liquid which results from the grinding operation into a subterranean earth formation zone.\nCertain problems arise, however, with respect to the disposal of certain types of wastes wherein the containment of the waste material in a particular portion of an earth formation zone is required. In the production of oil and gas, for example, a significant amount of naturally occurring radioactive material (NORM) may accumulate in the conduits and processing vessels associated with the production wells and related fluid handling facilities. The radioactive parts of these naturally occurring radioactive materials are usually Radium.sup.226 and/or Radium.sup.228. These radioactive substances are usually bound within the structure of scales and other produced solids. Typically, the composition of the scale is in the form of barium sulfate which is particularly hard and cannot be dissolved by most naturally occurring fluids. Accordingly, this barium sulfate scale, as well as other materials which are mixed with the scale, presents a disposal problem because of its radioactivity.\nCertain techniques for disposing of solids-laden slurries in subterranean formations are described in U.S. patent application Ser. No. 08/032,951 filed Mar. 17, 1993, by Thomas K. Perkins, et. al now U.S. Pat. No. 5,314,265, issued May 24, 1995, and assigned to the assignee of the present invention. This application describes a method of disposing of slurries of relatively fine particles of solid wastes into disposal zones which have a lower in situ compressive stress than that of an overburden zone and an underburden zone adjacent to the disposal zone. The disposal zone preferably has a permeability greater than the overburden and underburden zones. This technique is advantageous in that it is common to find earth formations having one or more relatively permeable zones of substantial thickness which are located below a formation zone comprising overburden of relatively impermeable material and which may exhibit in situ compressive stresses generally greater than the more permeable zone.\nU.S. patent application Ser. No. 08/043,323 filed Apr. 6, 1993, by Joseph H. Schmidt, et. al now U.S. Pat. No. 5,387,737, issued Feb. 7, 1995, and assigned to the assignee of the present invention describes a method of disposing of particulate solids in a subterranean formation by injecting a quantity of substantially solids-free liquid into a predetermined disposal zone having a predetermined porosity such that the injection of the solids-free liquid disaggregates at least part of the disposal zone, and then injecting a slurry of particulate solids into the disaggregated part of the disposal zone to place the solids in the pore spaces of the disaggregated part of the disposal zone. Other techniques which are of interest but which are not suitable for disposing of toxic wastes in accordance with the present invention are disclosed in U.S. Pat. No. 3,292,693 issued Dec. 20, 1966 to G. A. Hill, et. al, which describes a technique of disposing of liquid wastes into a low pressure zone in an earth formation which is between higher pressure zones and wherein liquids in the higher pressure zones tend to migrate toward the low pressure zone.\nStill further, Australian Patent 233,971, published Mar. 3, 1960 by J. J. Reynolds, et al, describes a method of underground disposal of radioactive liquids or slurries wherein a substantially horizontal fracture is created in a shale or nominally non-porous rock formation by opening the fracture with a suitable fracturing fluid and then injecting a radioactive waste material and propagating the fracture by injecting the radioactive fluid at a pressure greater than the formation breakdown pressure. Although generally horizontal fractures were thought to be the type of fracture typically formed in a subterranean earth formation according to the state of the art in 1960, it is generally accepted presently that fractures tend to extend in generally vertical planes under the urging of hydraulic or similar fluid injections. Moreover, totally non-porous or non-permeable formations may exist but they are difficult to locate and are usually not readily accessible for disposal operations.\nNotwithstanding the foregoing methods for disposing of slurries of waste materials, the present invention provides a method for disposing of toxic wastes, including radioactive materials, in an improved way which substantially satisfies the requirement of containment of the waste in a predetermined earth formation zone."}
{"text": "On-demand delivery of audiovisual content or other media content over broadband services may involve delivery of very large media files over broadband connections of varying speed and stability. Because available bandwidth may vary with network usage, location of the receiving device relative to a host or a server, and other factors, a delivery method that may be appropriate for one user at one location and at one particular time may not be satisfactory for another user at a different location or at a different time.\nFor example, while a first user at a first location may be able to stream motion video media content at high resolution at any time of day or night without playback stoppages, a second user at a second location may be able to satisfactorily stream motion video content only during low usage hours. During higher usage hours, the second user may instead download media content to avoid frequent playback stoppages. Progressive downloading, in which metadata about the media file is stored in the first part of the file to be transmitted, allows playback to begin before all of the media content is fully downloaded. Like streaming, however, progressive downloading permits the start of playback usually only after a certain portion of the media content has already been received in hopes of avoiding playback stoppages.\nAn additional problem arises when, during the delivery of the media file via streaming or downloading, the condition of the communication connection changes. For example, at a location where streaming was available at one time, increasing network usage or other changes affecting the communication connection may prevent satisfactory streaming. Conversely, it is possible that at a second location and at a second time, satisfactory streaming may not have been available but may become available."}
{"text": "This invention relates generally to an implantable containment apparatus made of selectively permeable material. In particular, the implantable containment apparatus is used to contain a therapeutical device, such as a drug delivery device, a cell encapsulation device, or a gene therapy device. A therapeutical device can be easily placed and replaced in an apparatus of the present invention without damaging tissues associated with the selectively permeable material of the apparatus.\nVarious implantable therapeutical devices, such as drug delivery, gene therapy, and cell encapsulation devices, have been disclosed over the years. A common feature of most of these devices is the use of selectively permeable, or semi-permeable, membranes to construct all or part of the device. These membranes contain their respective therapeutic agents and delivery systems within the particular device while being permeable to the desired therapeutical product. For cell encapsulation devices, the membranes are also permeable to life sustaining substances and to cellular waste products.\nWhen implanted in a recipient, the typical biological response by the recipient to most of these therapeutical devices is the formation of a fibrotic capsule around the device. With most drug delivery and gene therapy devices, this can limit the performance of the device, particularly when the therapeutic agent has a short half-life. For cell encapsulation devices, a fibrotic capsule encasing the device most often deprives the encapsulated cells of life sustaining exchange of nutrients and waste products with tissues of a recipient. The result is usually fatal to the encapsulated cells. Furthermore, a fibrotic capsule encasing a therapeutical device usually makes surgical retrieval of the device difficult.\nWhen certain therapeutical devices are implanted in a recipient, predominantly vascular tissues of the recipient can be stimulated to grow into direct, or near direct, contact with the device. On one hand, this is desirable because the therapeutical product of the device can then be delivered directly to the circulation of the recipient through the vascular tissues that are in contact with the device. On the other hand, this is undesirable because once vascular tissues of a recipient have grown in contact with one of these implantable therapeutical devices, removal of the device requires surgical dissection of the tissues to expose and remove the device. Surgical dissection of vascular tissues, particularly capillary tissue, can often be a difficult and painful procedure. Whether encased in a fibrotic capsule or surrounded with vascular tissue, the problem of retrieving these implanted devices is a considerable drawback of the devices.\nFor cell encapsulation devices, an alternative to retrieving and replacing the entire device in a recipient is to retrieve and replace the cells contained in the device. U.S. Pat. No. 5,387,237, issued to Fournier et al., is a representative example of a cell encapsulation device that has at least one opening into the device through which cells can be introduced and removed. Cells are introduced and removed in this, and other similar devices, as a suspension or slurry. Since most cell encapsulation devices are intended to correct a metabolite deficiency in a recipient caused by dysfunction or failure of certain of the recipient\"\"s cells, tissues, or organs, the source of the replacement cells is rarely the recipient. In a situation where non-autologous cells are used in this type of cell encapsulation device, the problem of contaminating a recipient with the foreign cells during loading, removal, or refilling of the device is ever present. One solution to this contamination problem would be to enclose the cells in a container that can be placed, removed, and replaced in a device as a unit.\nA retrievable cell encapsulation envelope enclosed in an implantable permselective membrane for use as an artificial endocrine gland is disclosed in U.S. Pat. No. 4,378,016 issued to Loeb. The Loeb device comprises a housing made of an impermeable hollow stem and a permselective membrane sack. The hollow stem has a distal end defining an extracorporeal segment, a percutaneous segment in the mid-region, and a proximal end defining a subcutaneous segment. The sack is adapted to receive an envelope containing hormone-producing cells and has an access opening that is coupled to the proximal end of the hollow stem. In a preferred embodiment, the cell containing envelope is in the form of a flexible collar. The flexible collar is partially collapsible to allow for easier placement and replacement of the envelope in the sack. Once in place, the flexible collar also provides a snug fit between the envelope and the sack. Placement and replacement of a cell containing envelope in the sack portion is accomplished manually with forceps, or the like. Retrieval of the envelope from the sack can be aided with a guidewire attached to the envelope. In one embodiment of the Loeb device, the sack has openings at both ends that are implanted percutaneously. In this embodiment, the cell containing envelope may be inserted or removed through either end of the device.\nThe housing of the Loeb device is surgically implanted in a recipient through the abdominal wall so the distal end of the stem protrudes from the recipient, the proximal end of the stem resides subcutaneously with respect to the abdominal wall, and the sack portion is placed in the peritoneal cavity surrounded in peritoneal fluid. According to Loeb, the sack allows hormones, nutrients, oxygen, and waste products to flow in and out of the sack while preventing bacteria from entering the patient. The sack and the envelope are said by Loeb to be permeable to nutrients and hormones, but impermeable to the hormone-producing cells and immune response bodies. Upon implantation of the device in a patient, the cells contained therein are said to take over the function of the corresponding natural gland, sense the amount of hormone needed, and produce the correct amount of the desired hormone.\nImplanted cell encapsulation devices, particularly those intended as an artificial endocrine gland, usually require a high rate of flux of nutrients and waste products between the encapsulated cells in the device and tissues of the recipient. Having a cell encapsulation device in close, or direct, association with a vascular structure usually provides the highest rate of nutrient and waste product flux for such a device. Loeb does not teach the value of vascularization of the sack portion of the housing, however. Nor is the Loeb device implanted in a part of the body that is particularly vascularized.\nBrauker et al. disclose a cell encapsulation device in U.S. Pat. No. 5,314,471 that requires close association of host vascular structures with the device. According to Brauker et al., xe2x80x9cconventional implant assemblies and methodologies usually fail to keep the implanted cells alive long enough to provide the intended therapeutic benefit.xe2x80x9d Cell death in these implanted devices is said by Brauker et al. to be due in large part to an ischemia imposed on the cells during the first two weeks following implantation. Brauker et al. conclude that xe2x80x9cthe cells die because conventional implant assemblies and methodologies themselves lack the innate capacity to support the implanted cells\"\" ongoing life processes during the critical ischemic period, when the host\"\"s vascular structures are not nearby.xe2x80x9d Brauker et al. state that in order for implanted cells to survive and function on a long term basis, the host must grow new vascular structures in association with the device. Brauker et al. note that a host will not naturally provide new vascular structures to an implanted cell encapsulation device. According to Brauker et al., the host must be stimulated by the implant assembly itself to grow new vascular structures close to the cell encapsulation device. Angiogenic stimulus can be provided by angiogenic factors applied to the cell boundary of the Brauker et al. device or by certain cell types encapsulated in the device. Growth of vascular tissue in association with a device of Brauker et al. will make removal of the device from a host difficult, however,\nAn implantable containment apparatus made of selectively permeable polymeric material that permits a therapeutical device, such as a drug delivery, a gene therapy, or a cell encapsulation device, to be placed and replaced in a recipient without damaging or disturbing tissues associated with the selectively permeable polymeric material would be useful. Such an apparatus that becomes closely associated with vascular structures without the need to supply angiogenic factors to induce the close vascularization would also be useful. A method of easily placing and replacing a therapeutical device in an implantable containment apparatus of the present invention would be further useful.\nThe present invention is directed to an implantable containment apparatus for a therapeutical device, such as a drug delivery device, a cell encapsulation device, or a gene therapy device. The apparatus is primarily made of selectively permeable material. The selectively permeable material permits flux, or exchange, of solutes between a therapeutical device contained in the apparatus and tissues of a recipient, while excluding cells from growing beyond a desired point through the material. When an apparatus is implanted in a recipient, various tissues from the recipient grow to associate with the apparatus. These tissues grow next to, or partially through, the exterior surface of the apparatus. It is preferred that vascular tissue is the predominant tissue that grows to associate with an apparatus of the present invention. Once growth of a recipient\"\"s tissues in association with an apparatus of the present invention has occurred, a therapeutical device can then be easily placed and replaced in the apparatus without damaging or disturbing the tissues associated with the selectively permeable material of the apparatus.\nThis is preferably accomplished by providing an implantable containment apparatus for a therapeutical device comprising a selectively permeable microporous polymeric material in the form of a tube wherein the tube comprises an exterior surface, an interior surface that defines a luminal space of substantially uniform diameter, and an access means at one end of the tube through which a therapeutical device is insertable into the luminal space of the tube wherein, once a therapeutical device is inserted into the luminal space of the tube, the therapeutical device is retained within the tube wherein biochemical and therapeutic substances having a molecular weight up to about 5,000,000 MW diffuse across the thickness of the tube between the contents of the therapeutical device contained therein and tissues of a recipient wherein the therapeutical device is removable from the tube through the access means of the tube and wherein the tube is refillable with a therapeutical device through the access means of the tube. The apparatus can also have an access means at each end of the tube. In this embodiment, a therapeutical device can be inserted and removed from the luminal space of the tube through either access means in the tube. Furthermore, in this embodiment, a fluid stream can be established through both access means of the tube that can then be used to flush a therapeutical device in and out of the luminal space of the tube.\nAccordingly, the present invention is also directed to a method in which a therapeutical device, such as a drug delivery device, a gene therapy device, or a cell encapsulation device, is easily inserted into, removed from, and replaced in a containment apparatus of the present invention as a unit with a fluid stream. The method involves repeatedly filling and emptying an implantable containment apparatus tube with a therapeutical device.\nOther features and advantages of the invention will become apparent upon review of the following specification, drawings, and claims."}
{"text": "1. Field of the Invention\nThis invention relates generally to bed structures for hauling; and more particularly, to a light-weight structure having improved clamp-down and tie-down mechanisms adjustably positioned thereon.\n2. State of the Prior Art\nVarious structures for supporting and transporting snowmobiles, all terrain vehicles, garden tractors, motorcycles, and the like have been described. Many of these structures are associated with trailers that can be towed behind a vehicle. Many general hauling structures included fixed tie-down devices around the periphery of the bed that are utilized in conjunction with chains, ropes, elastic style hooked tie-down devices, or retractable web tie-down material. Tying objects or various types of recreational vehicles in this manner is time consuming, inconvenient, and in many instances unsafe and inadequate.\nVarious types of specialty tie-down and clamp-down devices have been developed to be utilized in conjunction with specific objects. For example, various types of snowmobile trailers have bar structures that can be affixed across the tops of the snowmobile skis and are affixed to the edge frame of the trailer structure. These devices characteristically are in fixed position and do not accommodate a snowmobile being at various locations along the length of the bed. In some cases portions of the frame are adapted to receive part of skis and retain them. These arrangements are lacking in versatility, and in some cases do no adequately restrain the vehicles.\nCertain types of trailers utilized for transporting motorcycles and the like have a longitudinal channel in which the wheels are placed to prevent the wheels from sliding sideways, and the tie-down is accomplished by various devices as mentioned above. While assisting the positioning of the motorcycle, these channels substantially impair the use of the trailer for any other function, and do nothing to prevent the forward and backward motion of the motorcycle.\nMany of the transport beds are constructed of a supporting frame to which a deck or bed is affixed. Various structures and devices have been utilized to affix the bed to the frame, including bolts, clamps, screws and sundry fixtures and fasteners. The prior art type of fastening is costly in assembly and does not readily allow replacement of worn or damaged bed material. An improvement has been recognized in the prior art wherein side structural members have channels in which the bed material can be inserted and supported along its length. These prior art arrangements generally only deal with the entire bed and require substantial additional fastening structures across the width of the bed to the various frame support members. Such structures are not readily repairable or replaceable even with the sliding channels at the sides.\nIt is a primary objective of this invention to provide an improved transport bed structure.\nYet another objective is to provide a transport bed that has improved clamp-down and tie-down mechanisms.\nStill a further object of the invention is to provide a rotatable telescoping clamp-down mechanism.\nAnother objective of the invention is to provide an improved clampdown mechanism operable in conjunction with a retention channel for variable placement.\nStill a further object is to provide an improved tie-down mechanism operable with a retention channel.\nA further object of the invention is to provide an improved bed mounting structure.\nStill a further object of the invention is to provide bed mounting structures that minimize the cost of construction and maximize the ease of repair and replacement.\nAnother object of the invention is to provide improved clamp-down and tie-down mechanisms that are light-weight, durable, and relatively maintenance free.\nThese and other more detailed and specific objectives will become apparent from a consideration of the following description of the invention.\nThe present invention includes a bed structure utilized for transporting objects, and includes an improved clamp-down mechanism. One aspect of the improved clamp-down mechanism involves a telescoping clamping member that is associated with an elongated retention channel integrated as part of the bed structure, where the telescoping clamp mechanism can be selectively positioned along the length of the retention channel and utilized to clamp at various locations on the bed. One or more retention channels may be utilized in a spaced apart relationship. The telescoping clamp members can be retracted and rotated for installation into the clamp position and for releasing from the clamp position.\nA tie-down mechanism can also be used within the same retention channels and is likewise subject to variable positioning along the retention channel to form a tie-down relationship with an object.\nThe telescoping clamp-down mechanism and the tie-down mechanism are each retained in the retention channel by a rotatable cam that hold the mechanisms in the retention channel in a first position and allows them to be removed from the retention channel in a second position.\nThe retention channels and the structural frame members are adapted to have bed retaining channels along their lengths to support, retain, and protect the edges of the bed material without additional clamps or fastening mechanisms.\nThe rear frame member generally has a channel cross-section and is adapted for mounting lights within the channel recess to protect them from damage. In one configuration the rear member has a lower limit or lip portion for protecting the rear frame member if the associated bed structure is lowered to contact the surface for loading. The upper portion of the rear frame member includes a predetermined number of protrusions extending along the length thereof. A slidably engagable ramp structure is adapted with a longitudinal channel configured to match the arrangement of longitudinal protrusions on the rear frame member. The ramp member is slidably attached to the rear frame member when the bed structure has been put in place, and is adapted to retain and protect the rear edge of the bed. The ramp structure is sloped to allow ease of loading from the rear when the bed structure is utilized with a tiltable supporting structure.\nAll of the component parts are integrally formed from extruded light-weight material and cut to lengths required for assembly, there providing an efficient and cost-effective manufacturing system for various sizes of bed structures."}
{"text": "Generally, a communication application establishes a communication session between transmitting and receiving terminals connected to each other through a network and performs communications on the established session. If a propagation delay time between the transmitting and receiving terminals is very long or they communicate with each other across networks having different properties such as wired and wireless networks, then the throughput of communications between the transmitting and receiving terminals is lowered.\nOne known method of solving the above problem has been not to perform communications between transmitting and receiving terminals on one session, but to install a relay device between transmitting and receiving terminals and relay data between two sessions, i.e., a session from the transmitting terminal to the relay device and a session from the relay device to the receiving terminal. Examples of such a session relay process include “Indirect TCP” described in Non-patent document 1 (Ajay Bakre and B. R. Badrinath “I-TCP; Indirect TCP for Mobile Host”, Department of Computer Science Rutgers University, DSC-TR-314, 1994) and communication methods described in Patent document 1 (Japanese patent No. 3448481), Patent document 2 (Japanese patent No. 3482091), and Patent document 3 (Japanese patent No. 3494610).\nAn example of session communications according to TCP will be described below. Details of an ordinary TCP process are described in Non-patent document 2 (Jon Postel, “Transmission Control Protocol”, IETF, RFC 793, 1981) and Non-patent document 3 (W. Richard Stevens, “TCP/IP Illustrated, Volume 1: The Protocols, Addison-Wesley”, 1994, ISBN 0-201-63346-989).\nFIG. 1 shows a double-sided-transmission session relay device for relaying sessions transmissively from both a transmitting terminal and a receiving terminal. As shown in FIG. 1, double-sided-transmission session relay device 32-0 comprises packet input section 32-1 to which packets from a network are input, packet output section 32-2 which outputs packets to the network, session relay section 32-3 for terminating and relaying sessions, session status storing section 32-4 for storing a session status of session relay section 32-3, session determining section 32-5 for determining whether a packet is a session packet or not, session relay determining section 32-6 for determining whether a session relay process is to be performed or not, and session starting process monitoring section 32-7 for monitoring a session starting packet.\nPacket input section 32-1 performs a process of receiving packets from the network. Session determining section 32-5 determines whether an input packet from packet input section 32-1 is a session packet or not. If the input packet is a session packet, then session determining section 32-5 passes the input packet to session relay determining section 32-6. If the input packet is a packet other than a session packet, then session determining section 32-5 passes the input packet to packet output section 32-2.\nSession relay determining section 32-6 determines whether the session packet from session determining section 32-5 is a packet in a session registered in session status storing section 32-4 and also a packet sent after a session start acknowledging packet, or not. If these conditions are satisfied, then session relay determining section 32-6 passes the session packet from session determining section 32-5 to session relay section 32-3. If the conditions are not satisfied, then session relay determining section 32-6 passes the session packet from session determining section 32-5 to session starting process monitoring section 32-7.\nWhen session starting process monitoring section 32-7 receives a session starting packet, it registers session information representing a reception sequence number (a sequence number to be received from a receiving terminal) and a transmission sequence number (a sequence number to be transmitted to a transmitting terminal) which are temporarily generated, but not yet determined. When session starting process monitoring section 32-7 receives a session start acknowledging packet, it determines the reception sequence number and the transmission sequence number, which have been temporarily generated, in the session information. After the processing of session starting process monitoring section 32-7 is over, the session packet from session determining section 32-5 is passed to packet output section 32-2.\nSession relay section 32-3 performs a session relay process based on the session status stored in session status storing section 32-4, stores an updated session status in session status storing section 32-4, and passes a packet from session relay determining section 32-6 to packet output section 32-2, Packet output section 32-2 outputs the supplied packet to the network.\nA session according to a TCP in the session relay device shown in FIG. 1 will be described below. In the session according to the TCP, a session starting packet and a session start acknowledging packet are referred to as a SYN packet and a SYN• ACK packet, respectively.\nFIG. 2 is a block diagram of a network configuration including session relay device 32-0 shown in FIG. 1. FIG. 3 is a sequence diagram of a TCP relay process performed by session relay device 32-0 in the network configuration shown in FIG. 2. The example shown in FIG. 3 represents a sequence operation for transferring data from transmitting terminal 10 having IP address A to port number 80 of receiving terminal 20 having IP address B.\nFirst, transmitting terminal 10 sends a connection starting SYN packet to receiving terminal 20. The connection starting SYN packet includes information representing source IP address “A”, source port number “x”, sequence number “1”, destination IP address “B”, and destination port number “80”.\nIn session relay device 32-0 on the route, the connection starting SYN packet from transmitting terminal 10 is supplied to packet input section 32-1. Packet input section 32-1 passes the supplied SYN packet to session determining section 32-5. Session determining section 32-5 determines that the packet from packet input section 32-1 is a session packet, and passes the packet to session relay determining section 32-6. Session relay determining section 32-6 determines that the packet received from session determining section 32-5 does not corresponds to a session registered in session status storing section 32-4, and passes the received packet to session starting process monitoring section 32-7.\nSession starting process monitoring section 32-7 determines that the packet received from session relay determining section 32-6 is a session starting (SYN) packet, registers the session in session status storing section 1-4, and starts a session relay process to pass the received packet to packet output section 32-2. At this stage, at least a transmission sequence number to receiving terminal 20 and a reception sequence number from transmitting terminal 10 are stored in session status storing section 32-4. These two sequence numbers are initialized by the sequence number of the session starting (SYN) packet. The information of the session registered at this stage is incomplete session information as a whole because the transmission sequence number to receiving terminal 20 and the reception sequence number from transmitting terminal 10 have not yet been determined.\nHaving received the packet from session starting process monitoring section 32-7, packet output section 1-2 sends the received packet over the network to receiving terminal 20. The packet information in the SYN packet which is input to session relay device 32-0 and the packet information in the SYN packet which is output from session relay device 32-0 do not differ from each other.\nWhen receiving terminal 20 receives the SYN packet from session relay device 32-0, it returns a SYN• ACK packet to transmitting terminal 10.\nIn session relay device 32-0 on the route, the SYN• ACK packet from receiving terminal 20 is supplied to packet input section 32-1. Packet input section 32-1 passes the SYN• ACK packet to session determining section 32-5. Session determining section 32-5 determines that the packet received from packet input section 32-1 is a session packet, and passes the received packet to session relay determining section 32-6. Session relay determining section 32-6 determines that the packet received from session determining section 32-5 corresponds to a session registered in session status storing section 1-4, and is not a packet after the session start acknowledging (SYN• ACK) packet, and passes the received packet to session starting process monitoring section 32-7.\nSession starting process monitoring section 32-7 determines that the packet received from session relay determining section 32-6 is not a session starting (SYN) packet, but is a session start acknowledging (SYN• ACK) packet, updates the session information in session status storing section 32-4, and starts a session relay process to pass the received packet to packet output section 1-2. At this stage, the transmission sequence number to receiving terminal 20 and the reception sequence number from transmitting terminal 10 are updated. These two sequence numbers are initialized by the sequence number of the session start acknowledging (SYN• ACK) packet.\nHaving received the packet from session starting process monitoring section 32-7, packet output section 32-2 sends the received packet over the network to transmitting terminal 10. The packet information in the SYN• ACK packet which is input to session relay device 32-0 and the packet information in the SYN packet which is output from session relay device 32-0 do not differ from each other.\nThereafter, in session relay device 32-0, session relay section 32-3 returns an ACK packet responsive to the SYN• ACK packet to receiving terminal 20. In this manner, session relay device 32-0 opens a session between transmitting terminal 10 and receiving terminal 20.\nWhen the communication route suffers a certain failure after the session has been opened, a status notification is performed according to ICMP (Internet control massage Protocol). Details of operation according to ICMP are described in Non-patent document 3 and Non-patent document 4 (Jon Postel, “INTERNET CONTROL MESSAGE PROTOCOL”, IETF, RFC 792, 1981).\nIf an IP header contains a “Don't Fragment” flag when the route includes zones of different MTUs (Maximum Transfer Units), then the data cannot be fragmented and transferred. One solution to this problem is a process called Path MTU discovery proposed in Non-patent document 5 (J. Mogul, S. Deering, “Path MTU discovery”, IETF, RFC 1191, 1990).\nFIG. 4 shows a packet configuration of an ICMP Destination Unreachable/fragmentation request message. FIG. 5 shows a configuration of a PET part which has caused an error that is included in the ICMP Destination Unreachable/fragmentation request message. According to the Path MTU discovery process, if data cannot be fragmented and transferred because of the “Don't Fragment” flag in a router which includes zones of different MTUs, then the next-hop MTU is added to the ICMP Destination Unreachable/fragmentation request message shown in FIG. 4 and the message is sent to the source, so that the terminal of the source will retransfer a segment with a reduced MSS (Maximum Segment Size). In FIG. 4, “Internet Header+64 bits of Original Data Datagram” includes the information shown in FIG. 5 which is included in the packet which has caused the error.\nWhen the communication route contains a plurality of zones of different MTUs, a process of notifying the source of an MTU according to the ICMP to reduce the MSS is repeated until a minimum MTU is detected. At the time the minimum MTU is detected, communications are made possible between the source and the destination. A similar process is disclosed in Patent document 4 (JP-A No. 2003-18216).\nFIG. 6 is a diagram illustrative of an operation based on Path MTU discovery of a network apparatus having the double-sided-transmission session relay device shown in FIGS. 1 and 2. The operation based on Path MTU discovery with respect to an example according to TCP will be described below with reference to FIG. 6. In this example, the sum of IP and TCP header sizes is 40 bytes.\nAfter the SYN process, transmitting terminal 421 (address: A, port: x) sends data having a segment size of 1460 bytes and a packet size of 1500 bytes toward receiving terminal 423 (address: B, part: 80). The data output from transmitting terminal 421 is supplied to router 422 (address: C). Since the next hop is in the zone of “MTU=500”, router 422 which has received the data from transmitting terminal 421 adds the next-hop MTU to the ICMP Destination Unreachable/fragmentation request message, and sends the message to transmitting terminal 421. The ICMP Destination Unreachable/fragmentation request message from router 422 is supplied to transmitting terminal 421.\nIn response to the fragmentation request message, transmitting terminal 421 resends data having a segment size of 460 bytes and a packet size of 500 bytes toward receiving terminal 423. As the packet size of the data resent from transmitting terminal 421 is 500 bytes, router 422 sends the data resent from transmitting terminal 421 to receiving terminal 423. As a result, the packet of the resent data reaches receiving terminal 423, allowing data to be exchanged continuously between transmitting terminal 421 and receiving terminal 423."}
{"text": "1. Field of the Invention\nEmbodiments of the invention generally relate to a traffic control device or marker (hereinafter collectively referred to as a “traffic control marker”). More particularly, various embodiments of the invention are directed to a traffic control marker having a mesh base and an improved inner core.\n2. Description of the Related Art\nTraffic control markers used on roadways or other marking areas are frequently struck by moving vehicles. These traffic control markers typically have an outer or primary tube mounted to a base. Typically, the traffic control marker includes a reflective sheeting partially or completely surrounding the primary tube to warn or guide an operator of the moving vehicle, for example, at night or through a construction zone. Examples of commonly used traffic control markers include traffic cones and barrels, as non-limiting examples.\nAnother type of traffic control marker is a flexible delineator highway marker. The highway marker typically includes a flexible insert positioned in a primary tube of the highway marker to provide it with a rebound effect. In particular, the insertion of the flexible insert into the primary tube of the highway marker allows the traffic control marker to return to a substantially upright position after being struck and deflected by a moving vehicle. In addition to other attachment methods, a conventional traffic control marker may be attached to the roadway using an epoxy or other adhesive.\nA vehicular strike may overcome the seal created by the epoxy at the base of the traffic control marker, causing the traffic control marker to detach from the roadway at, or near, the base. Further, the primary tube, and thus the reflective sheet or indicia attached thereto, may detach from the traffic control marker upon impact by the moving vehicle. It can be costly and time consuming to repair or replace the damaged traffic control marker.\nThus, it would be desirable to provide a traffic control marker, or flexible highway marker, with an improved epoxy attachment mechanism capable of securing the primary tube to the base of the traffic control marker, thereby enhancing the structural integrity of the outer tube to prevent it from detaching from the traffic control marker on impact by a moving vehicle, improving the functioning life of the traffic control marker, and maintaining the performance of known types of flexible highway markers, when vehicles deflect them on the roadway or other marking area."}
{"text": "Fentanyl is a g-opioid receptor agonist with analgesic potency approximately 80-100 times that of morphine. In clinical settings, fentanyl exerts its principal pharmacologic effects on the central nervous system. Its primary actions are analgesic and sedation.\nThe analgesic effects of fentanyl are related to the blood level of the drug. In general, the minimum effective concentration and the concentration at which toxicity occurs rise with increasing tolerance to any and all opioids. The rate of development of tolerance may vary widely among individuals. All opioid mu-receptor agonists, including fentanyl, produce dose dependent respiratory depression. The risk of respiratory depression is typically less in patients receiving chronic opioid therapy who develop tolerance to respiratory depression and other opioid effects. Serious or fatal respiratory depression can occur, even at recommended doses, in vulnerable individuals.\nOrally administered fentanyl is subject to first pass effect metabolism as upwards of 50% or more of orally administered fentanyl is not absorbed. Other forms of delivery such a parenteral, buccal, and transdermal have been utilized to decrease or avoid this first pass effect for fentanyl.\nFentanyl is currently available in injectable form, as a lozenge (e.g. Actiq® fentanyl citrate; Actiq is a registered trademark of Anesta, LLC), and as a transdermal patch (e.g. Duragesic® 25, 50, 75, and 100 μg of fentanyl per hour; Duragesic is a registered trademark of Johnson & Johnson Corporation). Duragesic® provides continuous systemic delivery of fentanyl for approximately 72 hours. Duragesic® is indicated in the management of chronic pain in patients requiring continuous opioid analgesia for pain that is not optimally managed with lesser means such as acetaminophen-opioid combinations, non-steroidal analgesics, or prn (as needed) dosing with short-acting opioids. Duragesic® is typically not suitable for patients experiencing acute pain due to the delay in absorption of the fentanyl through the patch, or postoperative pain because serious or life-threatening hypoventilation could result.\nActiq® is a solid formulation of fentanyl citrate, intended for oral transmucosal administration. Actiq® is a lozenge attached to a handle similar in shape to a lollipop. The handle is purportedly to allow the Actiq® unit to be removed from the mouth if signs of excessive opioid effects appear during administration. Actiq® is indicated for the management for breakthrough cancer pain in patients with malignancies who are already receiving and who are tolerant to opioid therapy for their underlying persistent cancer pain. Actiq® is contraindicated in the management of acute or postoperative pain.\nSublingual tablets and lozenges (e.g., Actiq®) which may be used for acute pain or breakthrough pain have certain disadvantages. A disadvantage, amongst others, is that after intake the active agent in these pharmaceutical agents must first be released and dispersed prior to being available for resorption in dissolved form. In addition, the absorption pharmacokinetics of fentanyl from Actiq® may vary depending on the fraction of the dose that is absorbed through the oral mucosa and the fraction swallowed. Further, certain lozenges may be in the form of a candy which require medical supervision and may be socially questionable.\nThere exists a need in art for a sublingual formulation including fentanyl, a pharmaceutically acceptable salt thereof, or derivative thereof, which is suitable for effective pain management."}
{"text": "1. Field of the Invention\nThe present invention relates to a viewfinder brightness control circuit for a camcorder. More specifically, the present invention relates to a circuit which automatically controls the brightness of the image produced by an electronic viewfinder by applying an iris output voltage, which voltage varies in response to the sensed illumination of the environment.\n2. Discussion of Related Art\nWhen an object is photographed by a camcorder, the image quality produced by the electronic viewfinder is greatly affected by the illumination, i.e., the brightness of the surrounding. Accordingly, when an object is photographed by the camcorder, the viewfinder needs to be adjusted to compensate for the ambient illumination so that the image quality at the viewfinder is acceptable.\nAs shown in FIG. 3, a camcorder 100 includes a lens group 102 directing light through an iris 104 to a target 108. It will be appreciated that target 108 can be, e.g., a charge-coupled device (CCD) or similar device. Signals produced by target 108 are provided to a video signal processor 112 via a preamp 110. An indication of the brightness level at the target is provided from processor 112 to an iris motor 106, which is connected to iris 104, via an amplifier 114. Additionally, the video signal produced by video signal processor 112 is provided to an electronic viewfinder (EVF) 118 via control circuitry 116. In an exemplary case, the control circuitry included in block 116 is as shown in FIG. 1.\nIn general, brightness at the viewfinder is designed to be manually controlled. As shown in FIG. 1, a conventional viewfinder brightness control circuit includes an amplifier 1 and a light controller 3.\nAmplifier 1 includes a PNP transistor Q1 whose base is connected to a video signal input terminal via capacitor C1, whose emitter is grounded via a resistor R4, and whose collector is connected to a power voltage (Vcc) input terminal via a resistor R3; a resistor R1 connected between the power voltage input terminal and the base of PNP transistor Q1; a resistor R2 connected between the base of PNP transistor Q1 and ground; and a capacitor C2 having one terminal connected to the collector of PNP transistor Q1 and the other terminal connected to EVF 118 via light controller 3.\nLight controller 3 includes a resistor R5 connected between the other terminal of capacitor C2 (hereinafter node A) and the power voltage input terminal, a variable resistance VR1 and a resistor R6 connected to ground via the common node of capacitor C2 and resistor R5, and a capacitor C3 connected between the wiper of variable resistor VR1 and ground. It will be noted that the wiper of variable resistor VR1 is also connected to node A between capacitor C2 and resistor R5.\nVoltage at node A varies depending on variable resistor VR1 in light controller 3 and the video signal level output to the EVF 118 is determined in accordance with the variation in node A voltage. Thus, electronic viewfinder brightness is controlled by a video signal level output to EVF 118. In other words, if the viewfinder is located in well illuminated area, the voltage level of the video signal output to the electronic viewfinder will be high. When the electronic viewfinder is located in a poorly illuminated area, the voltage level of a video signal output to a viewfinder will be low. Accordingly, viewfinder brightness can be adjusted by controlling the resistance value of variable resistor VR1. It will be noted that the resistance value has to be selected in response to voltage level corresponding to the brightness level such that the voltage level is controlled after the user confirms the viewfinder brightness.\nAs described above, in the conventional operating method, the user directly controls variable resistor VR1 attached to the camcorder in order to control electronic viewfinder brightness. When the user operates the camcorder without manually adjusting the variable resistor, e.g., if the user depends on a preset brightness value stored in the camcorder, or when brightness is controlled after a video signal is displayed on the electronic viewfinder, the resolution of a viewfinder image quality can be deteriorated. This is a source of great inconvenience to a camcorder user."}
{"text": "This invention relates to a mixing element for mixing two or more fluids and a motionless mixer using the same.\nMotionless mixers (trademark: Static Mixers) have no mechanical moving parts and are so designed that fluids are mixed as they are passed through passages in a tube having a helical blade therein. The motionless mixers of this type are used in various fields, such as chemical plants, food industry, environmental pollution prevention technology, etc. In a prior art motionless mixer, as shown in FIGS. 1 to 3, a plurality of helical blades 1 and 2 formed by twisting sheets left- or right-handed respectively at an angle of 180 degrees are arranged in a passage tube 3. The blades 1 and 2 are coupled and fixed so that the facing end edges of each two adjacent blades cross at right angles. Two fluids flow individually through regions in the passage tube 3 partitioned by the blades 1 or 2, and are divided and joined at the junction of the blades. Then, the fluids flow helically, and are divided and joined at the next junction. After such division, joining, and transfer are repeated, the two fluids are mixed thoroughly.\nConventionally, the helical blades 1 and 2 are manufactured by casting, and each two helical blades 1 and 2 are fixedly coupled by welding or brazing so that their facing end edges cross at right angles. Helical blades fixedly coupled are inserted into the hollow cylindrical passage tube 3. Thus, the manufacture of the motionless mixer is not easy, and the fluids may often stagnate at the junctions where the adhesive agent is swollen. Owing to working errors, moreover, narrow gaps are formed between the helical blades 1 or 2 and the passage tube 3. These gaps reduce the efficiency of fluid mixing. As shown in the sectional view of FIG. 4 taken along line IV--IV of FIG. 3, the surface of the blade 1 and the inner peripheral surface of the passage tube 3 cross at substantially right angles, so that acute-angled dead spaces 4 are formed at four corner portions of the inside regions of the passage tube 3 partitioned by the blade 1. The fluids may stagnate at the dead spaces 4. Moreover, the flatness of an end edge 1a of the blade 1 increases the flow resistance.\nFIG. 5 shows a motionless mixer with helically twisted fluid passage tubes. As shown in FIG. 6, the fluid passage of each of passage tubes 5a and 5b has a semicircular cross section. The pair of passage tubes 5a and 5b are twisted helically, and are coupled so that their boundary between passage tubes 5a and 5b crosses the boundary between their adjacent pair of similar passage tubes 6a and 6b at right angles at the longitudinal end portions. Since it is difficult to manufacture a pair of helically twisted tubes, the manufacturing cost of this motionless mixer is high."}
{"text": "The present invention generally relates to ignition systems for gasoline engines. More particularly, the invention relates to a discharge ignition apparatus having triggering circuitry that generates a triggering signal based upon stored timing data.\nIgnition circuits of relatively elaborate design have often been provided to advance the ignition spark as engine speed is increased. For example, the spark may be delayed at starting speeds until approximately peak compression of the engine\"\"s piston. At higher engine speeds, the spark is preferably advanced to occur before peak compression.\nU.S. Pat. No. 5,392,753 to Burson et al., incorporated herein by reference, shows an ignition circuit utilizing a microprocessor.\nIn one aspect, the present invention provides an ignition apparatus for use with an internal combustion engine to produce an electrical spark at a spark ignition device. The apparatus comprises a magnet assembly, including a pair of pole faces, operatively revolved along a circular path. A magnetically permeable core is mounted adjacent to the circular path and has at least two leg portions each including a respective end face. The leg portions of the magnetically permeable core are situated such that the pole faces pass proximate to the end faces during revolution of the magnet assembly. As a result, a time-varying magnetic flux is produced in the magnetically permeable core.\nThe ignition apparatus further includes a transformer having a primary coil and a secondary coil related by a predetermined step-up ratio. The secondary coil is electrically connected during operation to the spark ignition device. A spark generation circuit is operative to apply a primary voltage pulse to the primary coil responsive to a triggering signal. The primary voltage pulse produces a spark generating pulse in the secondary coil.\nIn addition, the ignition apparatus includes triggering circuitry having a processor with stored timing data. The processor is operative to receive a signal indicative of the rotational position of the magnet assembly and produce a triggering signal. In a first mode, the triggering signal occurs at a fixed angle relative to top dead center. In a second mode, the triggering signal occurs at a variable angle relative to top dead center based on the stored timing data.\nOther objects, features and aspects of the present invention are discussed in greater detail below."}
{"text": "Reliable and consistent statistics relating to network traffic flow are important diagnostic tools in the operation of the modern datacenter. For instance, several classes of network errors and even failures include forwarding errors, very low rate packet loss that isn't accounted for by the device that is actually losing the packet, which may be identified if accurate packets statistics are available. However, current packet (and/or byte) counting methodologies and/or devices often fail to provide an accurate and consistent assessment of packet flow across a network, which negatively impacts the ability of this type of statistical information to serve the desired diagnostic purposes."}
{"text": "Prior attempts to provide information technology (IT) system coverage information, and IT management system information focus on management and coverage systems that stand alone. Each system has its own method of figuring out which resources it manages. Lists of resources generated by such systems may be incomplete, as each system is not aware of resources they do not manage. In recent years there has been an increase in the number of policies with which to comply. Compliance with the policies is only being checked for resources each system is aware of, creating the risk that a resource unknown to a system will create a noncompliance issue."}
{"text": "1. Field of the Invention\nThe present invention relates generally to improvements in portable spas, and more particularly, pertains to new and improved water features for a portable spa wherein a waterfall is provided as part of the portable spa for the enjoyment and entertainment of the spa users.\n2. Description of Related Art\nIn the field of portable spas, manufacturers are continually trying to enhance the portable spa product by adding various entertainment features to it beyond the hot water and jet massaging functions. Manufacturers have added audio systems to a portable spa device. Some manufacturers have added water features, such as waterfalls, to portable spa devices. The problem associated with the inclusion of a waterfall in a portable spa is the amount of space required by the waterfall. The present invention overcomes this problem."}
{"text": "In a typical communications network a wireless device, communicates via a Radio Access Network (RAN) to one or more Core Networks (CNs). The communications network may also be referred to as e.g. a wireless communications network, a wireless communications system, a communications network, a communications system, a network or a system.\nThe wireless device may be a device by which a subscriber may access services offered by an operator's network and services outside operator's network to which the operator's radio access network and core network provide access, e.g. access to the Internet. The wireless device may be any device, mobile or stationary, enabled to communicate over a radio channel in the communications network, for instance but not limited to e.g. user equipment, mobile phone, smart phone, sensors, meters, vehicles, household appliances, medical appliances, media players, cameras, Machine to Machine (M2M) device or any type of consumer electronic, for instance but not limited to television, radio, lighting arrangements, tablet computer, laptop or Personal Computer (PC). The wireless device may be portable, pocket storable, hand held, computer comprised, or vehicle mounted devices, enabled to communicate voice and/or data, via the radio access network, with another entity, such as another wireless device or a server.\nThe radio access network covers a geographical area which is divided into cell areas, with each cell area being served by a base station. The base station may be called a Radio Base Station (RBS), evolved NodeB (eNB), NodeB, B node, Radio Network Controller (RNC), Base Station Controller (BSC), Base Transceiver Station (BTS), MCE, depending on the technology and terminology used. A cell is a geographical area where radio coverage is provided by the radio base station at a base station site. The base station communicates with the wireless device(s) within range of the base station.\nAccording to the 3GPP, Multimedia Broadcast Multicast Services (MBMS) “is a point-to-multipoint service in which data is transmitted from a single source entity to multiple recipients. Transmitting the same data to multiple recipients allows network resources to be shared.” MBMS offers two modes: broadcast mode and multicast mode. The MBMS architecture enables efficient usage of radio network and core network resources. evolved MBMS (eMBMS) may be described as the Evolved Packet Core (EPC) version of MBMS. The eMBMS evolution brings improved performance thanks to higher and more flexible LTE bit rates, single frequency network operations, and carrier configuration flexibility.\nIn MBMS, there are some network nodes or functional entities which are important. Multi-cell/multicast Coordination Entity (MCE) is a network node or functional entity which is responsible for allocation of time and frequency resources for MBMS transmission. The MCE may be co-loated with for example an eNB. Another network node is the MBMS-GW, which is the entry point for incoming broadcast/multicast data traffic. The MBMS-GW broadcasts data packets to all eNBs within an area. Broadcast Multicast-Service Centre (BM-SC) is a network node or functional entity which is necessary in order for a communications network to support MBMS. The BM-SC is in charge of providing service to the end user.\nSome of the reference points in MBMS are Sn, SGmb and Sm. Sn is the reference point for the control plane between MBMS-GW and the SGSN. SGmb is the reference point for the control plane between BM-SC and the MBMS-GW. Sm is the reference point for the control plane between the Mobility Management Entity (MME) and the MBMS-GW. M3 Application Protocol (M3AP) supports the M3 interface which is between the MCE and the MBMS GW. A reference point may also be referred to as an interface. Signaling between nodes is exchanged at a reference point.\nThe purpose of a MBMS Session Start procedure is to request the radio access network to notify wireless devices about an upcoming MBMS Session of a given MBMS Bearer Service and to establish a MBMS Radio Access Bearer (RAB) and MBMS signalling connection for this MBMS Session. The MBMS Session Start procedure is triggered by the core network. For example, the core network initiates the procedure by sending a MBMS Session Start request message to the RNC. The MBMS Session Start request message comprises different parameters. The RNC acts according to the received MBMS Session Start request message. The RNC sends a MBMS Session Start response message or a MBMS Session Start failure message to the core network, depending on the outcome of the procedure.\nAccording to 3GPP, the loss or corruption of the data stored in the aforementioned network entities for support of MBMS service(s) for each MBMS session context which are created by the MBMS Session Start procedure and updated by the MBMS Session Update procedure will seriously degrade the MBMS service(s) offered to mobile subscribers. It is therefore necessary to define procedures to limit the effects of such failure, and to restore the MBMS service with minimized impact to the mobile subscribers. Such restoration procedures are related to failure and/or restart of several types of network nodes and network paths/interfaces, such as e.g. MBMS-GW, MME, SGSN etc. The terms restoration and re-establishment are equivalent.\nA failure may be a failure to receive a particular message, failure of a hardware or software component of a network node. A failure may be full/complete or partial. After a node has been restarted, all its bearer contexts are deleted.\n3GPP has started a new project called eMBMS restoration procedures, where the objective of this project is to specify enhanced restoration procedures to explicitly define the EPS behavior and to enable restoration of the eMBMS service when possible in order to minimize the end-user service impact upon different kinds of failure over the control path. Examples of such failures are as follows:                MBMS-GW failure/restart        MME/SGSN failure/restart        MCE failure/restart        BM-SC failure/restart        Sm/Sn path failure        M3AP path failure        SGmb path failure        \nIt has been agreed that the principle of these restoration procedures is to try to re-establish the control path to allow a subsequent MBMS session update and/or MBMS session stop. In addition, such re-establishing may be performed before bringing down the MBMS sessions. But the specific procedures under various failure scenarios are still under discussion.\nThe system behaviour upon restoration of a MBMS session towards an alternative downstream node during a SGmb path failure remains undefined in current standards."}
{"text": "The present invention involves the ejection of ink drops by way of forming gas or vapor bubbles in a bubble forming liquid. This principle is generally described in U.S. Pat. No. 3,747,120 to Stemme. These devices have heater elements in thermal contact with ink that is disposed adjacent the nozzles, for heating the ink thereby forming gas bubbles in the ink. The gas bubbles generate pressures in the ink causing ink drops to be ejected through the nozzles.\nThe resistive heaters operate in an extremely harsh environment. They must heat and cool in rapid succession to form bubbles in the ejectable liquid, usually a water soluble ink. These conditions are highly conducive to the oxidation and corrosion of the heater material.\nDissolved oxygen in the ink can attack the heater surface and oxidise the heater material. In extreme circumstances, the heaters ‘burn out’ whereby complete oxidation of parts of the heater breaks the heating circuit.\nThe heater can also be eroded by ‘cavitation’ caused by the severe hydraulic forces associated with the surface tension of a collapsing bubble.\nTo protect against the effects of oxidation, corrosion and cavitation on the heater material, inkjet manufacturers use stacked protective layers, typically made from Si3N4, SiC and Ta. Because of the severe operating conditions, the protective layers need to be relatively thick. U.S. Pat. No. 6,786,575 to Anderson et al (assigned to Lexmark) is an example of this structure, and the heater material is ˜0.1 μm thick while the total thickness of the protective layers is at least 0.7 μm.\nTo form a vapor bubble in the bubble forming liquid, the heater (i.e. the heater material and the protective coatings) must be heated to the superheat limit of the liquid (˜300° C. for water). This requires a large amount of energy to be supplied to the heater. However, only a portion of this energy is used to vaporize ink. Most of the ‘excess’ energy must be dissipated by the printhead and or a cooling system. The heat from the excess energy of successive droplet ejections can not raise the steady state temperature of the ink above its boiling point and thereby cause unintentional bubbles. This limits the density of the nozzles on the printhead, the nozzle firing rate and usually necessitates an active cooling system. This in turn has an impact on the print resolution, the printhead size, the print speed and the manufacturing costs.\nAttempts to increase nozzle density and firing rate are hindered by limitations on thermal conduction out of the printhead integrated circuit (chip), which is currently the primary cooling mechanism of printheads on the market. Existing printheads on the market require a large heat sink to dissipate heat absorbed from the printhead IC.\nInkjet printheads can also suffer from a problem commonly referred to as ‘decap’. This term is defined below. During periods of inactivity, evaporation of the volatile component of the bubble forming liquid will occur at the liquid-air interface in the nozzle. This will decrease the concentration of the volatile component in the liquid near the heater and increase the viscosity of the liquid in the chamber. The decrease in concentration of the volatile component will result in the production of less vapor in the bubble, so the bubble impulse (pressure integrated over area and time) will be reduced: this will decrease the momentum of ink forced through the nozzle and the likelihood of drop break-off. The increase in viscosity will also decrease the momentum of ink forced through the nozzle and increase the critical wavelength for the Rayleigh Taylor instability governing drop break-off, decreasing the likelihood of drop break-off. If the nozzle is left idle for too long, these phenomena will result in a “decapped nozzle” i.e. a nozzle that is unable to eject the liquid in the chamber. The “decap time” refers to the maximum time a nozzle can remain unfired before evaporation will decap the nozzle."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for the production of high purity calcium thiosulfate by a metathesis reaction of ammonium thiosulfate and calcium hydroxide or calcium oxide.\n2. Description of the Prior Art\nAlthough the prior art contains many references to the production of calcium thiosulfate, at present, there are no known commercial processes for the manufacture of high purity calcium thiosulfate.\nJapanese Pat. No. 6,039, issued Feb. 22, 1973, discloses a process for the manufacture of calcium or magnesium thiosulfate by reacting calcium or magnesium sulfite with sulfur in alkaline solution. High yields are obtained only in the preparation of magnesium thiosulfate. This Japanese patent also discloses a metathesis reaction between calcium hydroxide and magnesium thiosulfate to produce calcium thiosulfate and insoluble magnesium hydroxide. However, the preparation of calcium thiosulfate solutions according to this Japanese patent requires a difficult filtration of a slurry containing magnesium hydroxide.\nIn Spanish Pat. 245,171, issued Nov. 15, 1958, a process is disclosed for the preparation of calcium thiosulfate by reacting a mixture of calcium chloride hexahydrate and calcium hydroxide with a concentrated solution of sodium thiosulfate at 30.degree. C. Not only are 0.77 tons of co-product sodium chloride produced for every ton of calcium thiosulfate produced but residual sodium chloride is an objectionable contaminant, e.g., when calcium thiosulfate is used as a non-corrosive hardening accelerator for concrete.\nHigh purity calcium thiosulfate hexahydrate (CaS.sub.2 O.sub.3.6H.sub.2 O) is produced by a metathesis reaction between a concentrated solution of calcium chloride hexahydrate and sodium thiosulfate pentahydrate, described in Monatsh, 80, 220-31 (1949). A stream of nitrogen is bubbled into the reaction mixture to exclude air and CO.sub.2 and the pH of the metathesis reaction solution is adjusted to about 8 to 9 by addition thereto of calcium oxide. However, the solid calcium thiosulfate so produced requires a purification step at temperatures of 0.degree. C or lower.\nIt is accordingly an object of the present invention to provide a method for the production of high purity calcium thiosulfate by a metathesis reaction wherein difficult processing steps and objectionable impurities are avoided and conversion rates are high.\nIt is another objection of the present invention to provide a process for the production of aqueous solutions of high purity calcium thiosulfate up to the saturation concentration under conditions which minimize the decomposition of the thiosulfate ion.\nIt is still further an object of the present invention to provide a method for the production of calcium thiosulfate wherein integration of said method with a process for the manufacture of the ammonium thiosulfate feedstock is feasible.\nOther objects and advantages of the present invention will become apparent from the description which follows."}
{"text": "1. Field of the Invention\nThis invention generally relates to telecommunications and, more particularly, to maintenance termination units for locating trouble in a telecommunications network.\n2. Description of the Related Art\nMaintenance termination units help locate trouble in a telecommunications network. When trouble is reported, that trouble may lie within the customer's premises and/or within the local loop serving the customer. Maintenance termination units have been used to help determine whether the trouble lies within the local loop or within the customer's premises.\nU.S. Pat. No. 4,529,847 to DeBalko (issued Jul. 16, 1985) shows a maintenance termination unit. This maintenance termination unit comprises a pair of “normally open” voltage switches connected between the ring and tip of the local loop (see Column 3, lines 30–33). When this maintenance termination unit is installed at the customer's premises, the maintenance termination unit produces a distinctive, periodic DC signal (see Column 4, lines 37–50). When a fault lies within the customer's premises, this distinctive, periodic DC signal is produced regardless of polarity between the ring and the tip (see Column 4, line 63–Column 5, line 1). If trouble lies within the telecommunications network, then the maintenance termination unit produces either a steady signal (a “hard” fault) or a periodic signal (a “light” fault) (see Column 5, lines 1–7). Because this maintenance termination unit utilizes voltage switches, these switches unfortunately cycle between an open position and a closed position when indicating the presence of a fault (see Column 4, lines 46–48). This periodic signal is often confused by test technicians and leads to an inaccurate diagnosis. This constant cycling, between an open position and a closed position, also leaves little time to conduct a test of the local loop.\nOther maintenance termination units have disadvantages. Many of these prior art maintenance termination units utilize custom, integrated circuitry. These integrated circuits are expensive and time-consuming to design and to fabricate. The expense of designing and fabricating these circuits, in fact, is often cost-prohibitive. Even when an integrated circuit is used, the harsh, ambient environment leads to many field failures. When the integrated circuit fails, the entire maintenance termination unit is then wastefully discarded.\nThere is, accordingly, a need for an improved maintenance termination unit that is easier to use, less expensive to design and to manufacture, tolerant of ambient conditions, and cheaper to repair or replace in the field."}
{"text": "This invention relates to an electronic component for surface mount technology (SMT), having a component body comprised of an electrically insulating material, and at least one conductor element mounted on said body, which conductor element is comprised of electrically conducting material and has a contact region which is exposed to the exterior of the component body, which contact region serves to connect the component to a printed circuit board (PCB) or the like.\nCustomarily, an electronic component is connected to a PCB by inserting the conductor elements of said component into bores provided in the PCB, and then soldering the conductor elements in said bores.\nTo simplify the automated mounting of electronic components onto PCBS, in recent years a soldering technology has been developed which is known as SMT (surface mount technology), as a replacement for the traditional bore soldering technology. With SMT, the contact regions of the conductor elements are placed against a surface region of the PCB which region is typically coated with a soldering paste comprising materials which can be caused to produce a soldered joint. The soldering paste is then heated by IR radiation, to produce a connection between the electronic component and the PCB. This technique is known as the \"reflow method\". Alternatively, a vapor phase soldering process may be used, wherein the entire electronic component is enclosed in a space filled with a solvent and steam; the heat of condensation of the steam on the cold electronic component brings about the solder joint. However, this technique is objectionable because it poses environmental problems. The reflow method will now be described, with reference to FIGS. 3a and 3b of the accompanying drawings.\nThe fixing of an electronic component in the form of a plug connector, to a PCB (230, in FIG. 3a) will be described, with reference to FIG. 3a of the drawings: The plug connector is comprised of a body 210 of insulating material having a plate-shaped bottom member 212 and a wall 214 extending perpendicular thereto. Two conductor elements 220 pass through the bottom member 212. Outside of the body 210 the conductor elements have exposed contact regions 222, and inside of body 210 the conductor elements have plug-like regions 224 which can be inserted in a contact socket.\nTo fix the plug connector to the PCB according to the known \"reflow method\", the contact regions 222 are applied to a region of the PCB bearing a soldering paste, and then said paste is melted by means of infrared (IR) radiation applied to the contact region in the directions indicated by arrows 226, and the melted solder is allowed to solidify. The result is a strong and durable connection between the PCB and the plug connector.\nTo avoid shadow formation, caused by neighboring components, and which interferes with the action of the IR irradiation on the contact region, it has been proposed to chamfer the body 210 at the transition between the plate-shaped bottom member 212 and the wall 214 extending perpendicularly thereto (see inclined surfaces 218 FIG. 3b). This state of the art is disclosed, for example, in German OS 42 03 605 A1. If the other geometric parameters are left unchanged, the introduction of the chamfer increases the maximum angle of incidence, .alpha., with respect to the PCB (230), of the IR radiation incident on the contact region 222 without shadow formation. The chamfer feature enables a higher density of electronic components to be mounted on the area of the PCB using the surface soldering technique, due to the fact that the maximum angle of incidence of irradiation is increased by the chamfers 218, which in turn reduces the minimum separation between components which is required to avoid formation of shadows by neighboring components.\nAlthough the mounting density of components on a PCB can be greatly increased by the feature, the maximum achievable angle of incidence .alpha. is still limited, and thus the maximum achievable mounting density is limited, e.g. by the constraints imposed generally by the geometry of the component, namely the extent of the wall 214. Further, a problem is presented with the customary arrangement of electronic components fixed to a PCB by means of surface soldering techniques, namely that the range of angles at which the IR radiation may be supplied to the contact region is small, and the greater the mounting density the smaller becomes said range of angles. For a given IR source, this reduces the IR power available for melting, and thereby increases the time required to achieve melting."}
{"text": "In facility producing metal sheet materials, in particular a steel-sheet process line, when making the conveying rolls rotate at a high speed to run steel sheet, slip and walk of the steel sheet, sticking of dust and buildup on the surfaces of the conveying rolls, and other phenomena occur. At the time of processing, if slip of the steel sheet occurs, slip flaws are formed on the surface of the steel sheet, and the surface quality is impaired. If walk occurs, the processing speed falls and the productivity drops.\nFurther, if dust sticks and buildup occurs on the surface of the conveying roll, the shape of the dust and the foreign matter is transferred to the surface of the steel sheet whereby the surface quality is impaired and the grade of the steel sheet dramatically falls. Not only that, touchup to remove the foreign matter stuck to the surfaces of the conveying rolls becomes necessary and the productivity drops.\nSuch a problem is particularly remarkable in hearth rolls in continuous annealing furnaces processing sheets in a high speed atmosphere.\nSlip and walk of the steel sheet occur due to insufficient frictional force between the steel sheet and rolls. Therefore, to prevent slip and walk, it is necessary to increase the surface roughness of the conveying rolls to raise the coefficient of friction and to secure frictional force to suppress rising of the steel sheet. On the other hand, sticking of dust on the surface of the conveying rolls is the phenomenon of iron powder, sludge, or other foreign matter on the surface of the steel sheet sticking to the surfaces of the conveying rolls, while buildup is the phenomenon of iron, manganese oxide, etc. on the surface of the steel sheet sticking and building up on the surfaces of the conveying rolls.\nTo prevent these, it is effective to suppress sticking and catching of the iron powder, sludge, and foreign matter on the surface of the steel sheet causing sticking of dust and iron and manganese oxide etc. forming the source of buildup. It is necessary to reduce the surface roughness of the conveying rolls. That is, if adjusting the surface roughness of the conveying rolls to achieve both the prevention of slip and walk of the steel sheet and suppression of sticking of dust and buildup, optimization becomes difficult due to the contradictory nature of the measures.\nAs measures for suppressing slip and walk of the steel sheet and sticking of dust and buildup on the surfaces of the rolls, the method of forming a flame sprayed coating on the surfaces of the rolls, then dulling it (Japanese Patent Publication (A) No. 7-57904), the method of dulling the surface, then forming a flame sprayed coating (Japanese Patent Publication (A) No. 7-47766 and Japanese Patent Publication (A) No. 9-157826), the method of forming a flame sprayed coating, then polishing it (Japanese Patent Publication (A) No. 10-168527 and Japanese Patent Publication (A) No. 10-168528), and the method of dulling the surface, then forming a flame sprayed coating and further polishing it (Japanese Patent Publication (A) No. 7-22773, Japanese Patent Publication (A) No. 7-1021, Japanese Patent Publication (A) No. 7-39918, and Japanese Patent Publication (A) No. 2005-105338) have been proposed.\nHowever, with each of these methods, achievement of both prevention of slip and walk of the steel sheet and prevention of sticking of dust and buildup on the surfaces of the rolls was difficult.\nThat is, in the method proposed in Japanese Patent Publication (A) No. 7-57904, lasers are used to drill holes. The side walls of the holes become sharp, so clogging by dust more easily occurs. Further, the lasered parts suffer from microcracks, so the surface roughness becomes larger and sticking of dust and buildup easily occur. Further, in the method proposed in Japanese Patent Publication (A) No. 7-47766 and Japanese Patent Publication (A) No. 9-157826, the surface of the flame sprayed coating becomes a rough surface as flame sprayed and sticking of dust and buildup easily occur. Further, with the method proposed in Japanese Patent Publication (A) No. 10-168527 and Japanese Patent Publication (A) No. 10-168528, the surface roughness Ra becomes smaller and slip and walk of the steel sheet easily occur. Further, the method proposed in Japanese Patent Publication (A) No. 7-22773, Japanese Patent Publication (A) No. 7-1021, and Japanese Patent Publication (A) No. 7-39918 reduces the surface roughness of only projecting parts formed at the surface of flame sprayed coatings. There was the problem that when conveying steel sheet softened at a high temperature, the inclined parts near the projections contacted the steel sheet and easily formed buildup at the rolls.\nFurther, Japanese Patent Publication (A) No. 2005-105338 proposes shot blasting or brushing to treat the surface of the flame sprayed coating and set the ratio R/R′ between a roughness parameter R (any of Ra, Rq, Rp, Rv, and Rz) measured based on JIS B0633 setting the cutoff value to an initial value and a roughness parameter R′ measured setting the cutoff value to 1/10 of the initial value to 4 or more so as to suppress buildup. However, with this method, the fine relief of the surface of the flame sprayed coating could be reduced, but there was the problem that the shot blasting or brushing ended up causing the formation of a new modified layer on the surface of the flame sprayed coating and due to this sticking of dust and buildup could not be completely prevented."}
{"text": "Distributed processor architectures are in use in base stations of radio systems, both at the level of the whole base station and at the unit level. In distributed processor architectures, several processors can attend to similar tasks, such as in digital signal processing, where several digital signal processors are connected to each other. Sub-units of such distributed processor architectures can also contain other components, for instance memories and application-specific integrated circuits (ASIC).\nUsually, such sub-units connected to each other boot automatically. Each sub-unit and a processor in it can be booted independently, for instance by means of boot software stored in a non-volatile memory, such as a flash memory, connected to each processor.\nAs regards sub-processors executing similar functions, their boot software is usually the same. The disadvantage of the prior art solutions is that non-volatile memories with identical boot codes must be connected to all sub-processors. Thus, more components are required and much printed board space is consumed. Therefore, also costs are increased. Further, when a printed board of a unit formed by the process architecture is manufactured, more solder joints are required, whereby also the manufacturing is slower and the sensitivity to failure and the failure density in the device increase."}
{"text": "In any communication environment including a central equipment and a plurality of remote equipments, especially in digital communication systems operating in multipoint, the central equipment has to control frequently the presence of each remote equipment. Indeed, in any hierarchical system, internal resources are allocated for each remote equipment. The management of such resources must be performed as rapidly as possible after use, as soon as the remote equipment disappears. Such resources can be control blocks, traffic queues, timers, semaphore software . . . As the total resources of the central equipment are size limited, the non-release of unused resources can result in refusing to serve new remote equipments.\nA classical method to avoid such a drawback consists for the central equipment to poll successively the remote equipments and invite them to transmit an information confirming they are alive. Unfortunately, the polling protocol results in an overhead, that is the portion of bandwidth consumed to implement the protocol, which is all the more important since the time required to detect that a remote equipment has disappeared is short."}
{"text": "Network switches (also known as a switching hubs, bridging hubs, MAC bridges or simply “switches”) are computer networking devices that connect other devices together in a computer network. Network switches are typically used in complex data networking environments to, for example, interconnect various devices within one or more subnetworks. Some network switches use packet switching to receive, process, and forward data to destination devices. Switches can be used for electrical switching, optical switching that enables signals in optical fibers or integrated optical circuits to be selectively switched from one circuit to another, or a combination of electrical and optical switching.\nTypical network switches include a plurality of circuit boards with associated switch circuitry (e.g. line cards and fabric cards) that are interconnected via backplane or midplane circuit boards within an enclosure. Although some cards include handles for manually removing the cards from the enclosure, conventional handles (particularly for optical cards) do not have locking features to hold the cards in place. Additionally, such handles typically require manual manipulation to eject the cards from the enclosure.\nAs is known, line cards typically include a modular electronic circuit designed to fit on a separate printed circuit board (PCB) and interface with a data communications network, and can provide transmitting/receiving ports for a local area network (LAN) and/or a wide area network (WAN). Fabric cards can include one or more switch elements for implementing the stages of the switch fabric. The line cards and fabric cards can be mounted to the backplane and/or midplane in a motherboard/daughterboard relationship in which the backplane or midplane extends transversely across the enclosure. In this arrangement, the backplane/midplane is perpendicular to the flow of cooling air through the enclosure, and can block or otherwise impede the flow of cooling air, resulting in relatively high operating temperatures that can reduce the longevity and/or performance of switch components. Another factor that can affect cooling is the efficiency of the air inlets and air outlets on the switch enclosure. Conventional network switches utilize fans to move cooling air through the enclosure and maintain the electronic components at suitable operating temperatures. Typically, perforated sheet metal is used to form air inlet and outlet vents on the front and back of the enclosure, respectively. Although such material can provide suitable covers for the enclosure, limitations on the size of the perforations limit the amount of air flow into the enclosure and can result in less than optimum cooling.\nIn conventional network switches, the central processing units (CPUs) are typically mounted to line cards. As a result, upgrading CPUs in a conventional switch generally requires removal and replacement of line cards to change out the CPUs. CPU technology, however, tends to advance at a relatively quick pace, which can lead to relatively frequent line card removal and replacement for CPU upgrades in conventional switches. In view of the shortcomings associated with conventional network switches, it would be advantageous to develop a network switch that provides, among other things, efficient cooling and relatively easy component replacement."}
{"text": "Security systems are presently limited in their ability to detect contraband, weapons, explosives, and other dangerous objects concealed under clothing. Metal detectors and chemical sniffers are commonly used for the detection of large metal objects and some kinds of explosives, however, a wide range of dangerous objects exist that cannot be detected with these devices. Plastic and ceramic weapons developed by modern technology increase the types of non-metallic objects that security personnel are required to detect. The alternative of manual searching of subjects is slow, inconvenient, and would not be well tolerated by the general public, especially as a standard procedure in, for example, airports.\nKnown prior art X-ray systems suggested for detecting objects concealed on persons have limitations in their design and method which prohibit them from achieving the low dose and high image quality which are prerequisites of commercial acceptance.\nRadiation exposure is an important consideration in X-ray concealed object detection systems. The United States National Council on Radiation Protection (NCRP), in NCRP Report No. 91, \"Recommendations on Limits for Exposure to Ionizing Radiation\", 1987, addresses this issue. In this report, the NCRP states that a radiation exposure of less than 1000 microRem per year in excess of environmental levels is negligible, and efforts are not warranted at reducing the level further. Persons employed in security facilities or frequently traveling by airlines may be subjected to many hundred security examinations per year. A yearly radiation exposure of 1000 microRem permits a single scan exposure within the range of 1 to 10 microRem for the general public. In accordance with the NCRP recommendations, radiation levels significantly higher than this present a non-trivial health risk.\nAn inspection system which operates at a low level of radiation exposure is limited in its precision by the small number of X-rays that can be directed against a person being searched. X-ray absorption and scattering further reduces the number of X-rays available to form an image of the person and any concealed objects. In prior art systems, this low number of detected X-rays has resulted in unacceptably poor image quality.\nRadiant energy imaging systems have been proposed to detect concealed objects. The system disclosed in U.S. Pat. No. Re. 28,544 of Stein et al. projects a scanning pencil beam of X-rays through the subject's body where the beam will be transmitted or absorbed depending upon the concealed object, if any. A detector is scanned vertically behind the subject to collect transmitted X-rays. Stein et al. comment that, \"although the detector 25 is shown behind the object being scanned for responding to the radiant energy transmitted through the object being scanned, it is within the principles of the invention to position the detector in the region between the radiant energy source and the object being scanned to respond to the scattered energy.\" However, the X-ray tube potential is set at up to 150 Kilovolts and is specifically chosen to transmit X-rays through the person being examined. Operation of the X-ray tube at 150 Kilovolts or even at 100 Kilovolts would negate benefits of imaging by backscatter detection, e.g., low dose scanning. This technique requires the subject to be exposed to a substantial radiation dosage, especially if the subject is scanned often, e.g., a frequent flyer.\nU.S. Pat. No. 4,839,913 of Annis et al. teaches a flying spot scanning system for baggage inspection which uses X-ray-induced fluorescence to permit detection of concealed objects. Since fluorescence is dependent upon atomic number, each object emits a fluorescent radiation line which is unique to its atomic number Z. Scattered, transmitted and fluorescence signals are generated to distinguish objects from the background. The energy level of the source must be sufficient to excite a selected fluorescence radiation line so that the fluorescent radiation line has sufficient energy to escape the object. Thus, it may be necessary to expose the object to relatively high X-ray energy in order to detect certain materials, which would be unacceptable for personnel inspection systems.\nAnother baggage inspection system is disclosed (U.S. Pat. No. 4,799,247 of Annis et al.) which has detectors for both transmitted and backscattered X-rays to independently produce signals from the same incident beam. The separate signal may then be used to enhance each other to increase the system's accuracy in recognizing low Z materials. Clearly, with the incident beam being of sufficient energy to provide both transmitted and backscattered signals, the X-ray energy must be relatively high, similar to that of Stein et al., making such a system undesirable for personnel inspection. A similar technique is described by Glockmann et al. (U.S. Pat. No. 4,884,289) for baggage inspection where both transmission and backscatter is used.\nSystems of the foregoing prior art that apply to body scanning are designed to detect radiant energy that has been transmitted through the body, scattered from the body, or emitted from the body. In the prior art systems, images are produced by characteristics of the subject's body and any object concealed under the subject's clothing. The system operator then inspects each image for evidence of concealed objects.\nFor various reasons the prior art systems do not adequately detect plastics, ceramics, explosives, illicit drugs, and other non-metallic objects. One reason in particular is that these materials share the property of a relatively low atomic number (low Z). Low Z materials present a special problem in personnel inspection because of the difficulty in distinguishing the low Z object from the background of the subject's body which also has low Z. Detection of low Z materials is addressed in the Annis '247 patent, but only for baggage. As mentioned above, this technique exposes a live subject to radiation levels which are excessive.\nIt would be desirable to provide an X-ray inspection system for detecting concealed objects carried by a person which is capable of detecting low Z materials as well as metals, but which does not expose the subject to radiation doses significantly higher than normal environmental radiation levels. It is to such a system that the present invention is directed."}
{"text": "Search engines use several techniques for sorting or sorting pages arising from a search. Among the known techniques for exploring a set of web pages, some rely on semantics, a page being ranked as being all the more relevant if it comprises a large number of occurrences of the word or words searched for. These techniques are sensitive to a practice, known as “spamming”, aimed at causing the words commonly used by surfers in their search query to appear a very large number of times in a given page, this having the effect of causing the page to appear frequently to be relevant.\nOther techniques are based on the topological structure of the web. These techniques take account at one and the same time of the existing links between the pages considered and properties of the pages themselves, such as whether a page belongs to a domain or to a network subdomain of the web. These techniques are generally based on a graph representation of the pages to be processed. They are appropriate to the classification of pages complying with given topological properties in the graph. These techniques are sensitive to a variant of the practice of “spamming” aimed at referencing a given page a large number of times, this having the effect of locally falsifying the topological characteristics of the graph of the web.\nSome of the techniques utilizing the topological structure of the web consist in effecting a classification of the web pages by allocating to the various pages a rank which is dependent on the relations of one page with the others.\nAn example of such a procedure, known by the term “PageRank”, is used in the implementation of the Google™ search engine and is described in the document: “The PageRank Citation Ranking: Bringing Order on the Web”, by L. Page, S. Brin, R. Motwani and T. Winograd; Technical Report, Computer Science Department, Stanford University, 1998.\nThe PageRank procedure orders the pages as a function of their visibility on the web. In this procedure, random navigation from page to page around the web following the hypertext links is simulated. This navigation corresponds to that operated by a user accessing the web when the latter randomly activates one of the hypertext links located in a displayed page, so as to access another page. This procedure carries out a probabilistic analysis of this simulated navigation so as to determine the probability of the user being located in a given page during random navigation from page to page such as this. The rank of a page is all the higher the higher the number of times this page is cited by other pages.\nSuch a procedure provides a ranking which is not necessarily relevant as regards the search performed by a user, the best ranked pages (of highest rank) not necessarily being the pages corresponding best to the user's expectation.\nFurthermore, this procedure does not make it possible to identify in the set of documents thematic communities or communities of interest, apt to steer the user more rapidly to an interesting page, nor even to perform a classification of the documents found by thematic community."}
{"text": "(a) Field of the Invention\nThe present invention relates to a heat sensitive transfer medium for use in thermographic recording which comprises combining a color former-containing transfer sheet with a developer-containing receiving sheet.\n(b) Description of the Prior Art\nAs the heat sensitive transfer medium for use in thermographic recording, there have hitherto be known the one which comprises a transfer sheet formed by providing a heat-sublimable dye layer on a substrate and a receiving sheet which receives a dye image sublimed and transferred from said transfer sheet by heat-printing from the back (the substrate surface) of the transfer sheet, and the one which comprises combining a transfer sheet formed by providing a heat-fusible material and pigment or dye-containing transfer layer on a substrate with a receiving sheet which receives an image fused and transferred from said transfer sheet by heat-printing from the back of the transfer sheet. The former transfer medium is defective in that due to the use of the heat-sublimable dye, the dye image on the receiving sheet is inferior in preservability and consequently there is a necessity of overcoating the transferred image, while the latter transfer medium is defective in that because of having the transfer layer formed by dispersing pigment or dye in the heat-fusible material, when a large amount of pigment is contained in said layer in order to obtain a high density image, the transfer efficiency deteriorates, whereby it becomes difficult to obtain a high density image, and when using a large amount of heat-fusible material in order to improve the heat sensitivity, furthermore, as the result of said large amount of heat-fusible material transferring on the side of the receiving sheet, the separation of the transfer sheet from the receiving sheet can not be done smoothly to cause the fine line image area to lack definition.\nOn the other hand, a heat sensitive transfer medium (for instance, such as the medium disclosed in British Pat. No. 1,053,905), which comprises supporting on a substrate separately two kinds of materials, for instance the color forming agent such as benzoindolinospyropyran and the developer such as phenolic compound, which form color by their mutual thermal reaction, and contacting these supporting layers (transfer layer and receiving layer) facing each other to effect heat-printing, is also known, but the medium of this type is observed to be defective in that since the medium of this kind is of reaction type, if the transfer layer transfers only toward the receiving layer at the time of facing contact, the color formation is not effected to the full, thereby producing a low density image, while if the heating conditions are effected at higher temperature and for longer times in order to promote the sufficient reaction, the image on the receiving sheet surely becomes a high density image, but the color forming reaction proceeds also on the transfer sheet to thereby cause image formation.\nWe previously proposed, as the heat sensitive transfer medium which had improved the above mentioned drawbacks, the one prepared by incorporating a large amount of oil-absorptive porous filler in the receiving layer formed on the receiving sheet. The heat sensitive transfer medium of this type surely can obtain a high density image by a large number of transfer, but is still not satisfactory in respect of the uniformity of image density.\nIn addition, it was difficult for the above mentioned conventional heat sensitive transfer mediums to obtain images which are superior in the mechanical strength such as friction resistance, abrasion resistance and the like."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus primarily for the purpose of nonintrusive scanning of cargo containers for nuclear based weapons which is intended for use during container transfer between a ship and land transportation.\nMore particularly it relates to a buffer crane having radiation emission scanning apparatus arranged for nonintrusive interrogation or inspection of cargo containers while each container is being transferred between a ship and quay side land transportation without slowing the quay crane container transfer cycle.\nStill more particularly, the present invention is a mobile cargo container buffer and scanning crane which transfers cargo containers between land transportation and a quay crane pickup buffer position and which suspends the cargo container being transferred at a predetermined inspection position for radiation emission examination by a craneboard apparatus before depositing the cargo container at the intended transfer buffer deposition position.\nSpecifically it relates to a mobile platform having a bridge crane mounted thereon for transferring cargo containers between either a quay crane pickup position, or quay side land transportation, and a predetermined inspection position on said platform where it can be non-intrusively inspected by a longitudinally reciprocating radiation emission inspection apparatus while functioning as a buffer operation. The container is then either deposited on land transportation such as a truck trailer chassis or deposited on the buffer position for pickup by a quay crane for transfer to a ship.\n2. Description of the Prior Art\nIn view of recent terrorist activities throughout the world, considerable effort is being given by analysts to improving security measures with respect to the maritime industry and United States port operations. At the present time, based on the Automated Tracking System, an intelligence based search system used by U.S. Customs, it has been estimated that the probability that contraband will be discovered in containers that are entering the U.S. is less than 50 percent. Thus, the maritime industry provides a delivery system for weapons of mass destruction and every other container could possibly conceal an atomic weapon which, if exploded in a U.S. port, in addition to causing massive destruction, would essentially end international trade.\nThe methods to be employed to improve U.S. port security may prove severely detrimental to port productivity. The ideal performance standard is 100 percent inspection of inbound containers to U.S. ports. Rather than physically unloading and inspecting the contents of every container, a more efficient alternative is to non-intrusively inspect each container such as by x-ray, and when more sophisticated means of nonintrusive inspection are developed, such as gamma ray scanning and neutron analysis, implementing the use of those developments in addition to x-ray based systems.\nThe use of x-ray machinery in one manner or another for the purpose of inspecting containers is discussed in the prior art literature. However, implementing prior scanning systems has essentially added processing steps to port operations, and the size of the machinery and complexity of the processing steps have essentially interrupted the established system of port operations. This factor causes problems when attempting to integrate the technology into cargo container handling.\nThe problem with instituting high-energy x-ray scanners for cargo container examination, apart from safety concerns, is that the equipment will reduce port productivity by disrupting highly developed port operations and consuming valuable terminal space. The primary disadvantage is the interruption of the cargo container transfer process between ship and shore. Most importantly, it interferes with the quay crane offloading cycle time which is crucial to a ship's berthing time at the dock which must be kept at a minimum.\nThe inspection procedure requires holding a container immobile so that it can be x-rayed. While the cycle times for the x-ray process may ultimately be reduced by improved technology, all of the presently considered means for effecting the x-ray process require either stopping the container movement for processing (usually during unloading in U.S. ports and, if required, during loading at foreign ports) or multiple additional handling steps of the container during the transfer process by taking it out of the normal handling cycle, and to an extra handling step, at an x-ray position for processing, and then returning the container into the transfer cycle.\nThe present invention permits integration of the x-ray process into the buffer station method of crane operation disclosed in applicant's above-referenced patents such that the inspection process can occur concurrently with the cargo container transfer between ship and shore without interruption of the quay crane transfer cycle.\nThe mobile cargo container scanning crane contemplated according to the present invention departs substantially from the conventional concepts and designs contemplated by the technical literature, and, in doing so, provides an apparatus primarily developed for the purpose of nonintrusive cargo container inspection during transfer between ship and shore as described above, but it accomplishes the result in a different and improved manner by producing a transfer cycle with a buffer-inclusive procedure for container inspection which is easily integrated into the container transfer cycle for faster processing times and more efficient port operation."}
{"text": "1. Field of the Invention\nThe present invention relates to the side airbag apparatus, which protects the occupant of a vehicle against injury in the event of side collision by inflating the airbag into a vehicle cabin. More specifically, the present invention relates to the side airbag apparatus, which can surely protect the occupant of a vehicle against injury by inflating airbags into a vehicle cabin from a seat of a vehicle.\n2. Description of Revelant Art\nThe airbag apparatus, which deploys the airbag instantly and protects the occupant of a vehicle against injury in the event of collision of a vehicle, has been well known. Among this kind of airbag apparatus, the side airbag apparatus, which deploys the airbag towards the region between the occupant and the door so as to protect an occupant of a vehicle against injury, has been used against the occurrence of the side collision.\nAs an example of these types of conventional side airbag apparatus, there is provided the side airbag apparatus disclosed in Japanese unexamined patent publications H9-323 607.\nIn the conventional side airbag apparatus, since the airbag is inflated by high-pressured gas, which is supplied at constant rate, the rate of inflation and the pressure of the airbag cannot be controlled. Thereby, in case of the airbag apparatus equipped with two airbags, if one airbag deploys toward a weak part, such as a flank and the other airbag deploys toward a strong part, such as a leg, of the body, the weak part of the body may be damaged because both airbags are inflated at the same rate of inflation."}
{"text": "This invention relates to motor vehicle accessories, and more particularly, to reflective safety devices."}
{"text": "A refining element is generally formed with a pattern of bars and intermediate grooves. The bars and grooves are formed in different ways, depending on which fibrous material is worked and the degree of refining desired and, thus, in the case of lignocellulosic material, the pulp quality which is desired. The bars have an upper surface and side surfaces so that longitudinal edges are formed between the upper surface and respective side surface. The bars can be, for example, continuous or discontinuous, and can be arranged in different patterns. The working of the fibrous material is substantially carried out by the bars of the refining elements. The refining gap is formed so that the fibrous material, seen in the radial direction, passes from the inside outwardly. Farthermost inward in the refining gap the refining elements are normally formed so as to bring about a first disintegration of the material and to advance the material outwardly in the refining gap. A certain defibering, i.e. separation of the fibers of the lignocellulosic material, also takes place in the inner portion of the refining gap where the distance between the refining surfaces is the greatest. Thereafter, that distance decreases outwardly in order that the desired working or refining of the fibrous material shall be obtained.\nDuring the refining of fibrous material of high concentration, and above all at high energy inputs, it has been found necessary to place flow restrictions, so-called dams, in the grooves of the refining elements in order to prevent unworked material from passing out through the refining gap. These dams, however, form an obstacle for the steam developing in the refining gap during the refining. A high steam pressure is thereby created in the refining gap. This high steam pressure has a negative effect on the capacity and operational stability of the refiner. It also implies a limitation on the possible energy input. The developed steam, thus, will be forced by the flow restrictions upwardly out of the grooves, and will disturb the material flow through the refining gap.\nOne way of solving this problem would be to supply dilution water to the refining gap in order to thereby condense the steam. This, however, results in reducing the material concentration to a low level, and thus in a deteriorated pulp quality.\nDuring the working or refining of fibrous material with low concentration no steam development takes place, and the material is partially transported by the liquid flow out of the refining gap. Also in this case dams are usually used to prevent unworked material from passing out through the refining gap. It can imply, however, that the flow through the refining gap will be much too low."}
{"text": "Visual inspection is widely used in the tire manufacturing process and even more commonly relies on the skill of the operators tasked with checking for the absence of visible imperfections at the surface of the tires in order to ensure the compliance thereof.\nHowever, with the advances in the computing power of computer-based means, tire manufacturers are developing automatic inspection means to assist the operators tasked with the visual inspection. To this end, it is possible to use an inspection device comprising lighting means and cameras which are positioned in such a way as to scan the exterior and interior zones of the lateral beads and of the tread of the tire that is to be inspected. The viewing field of each camera is angularly limited. In order to obtain complete images of the inside and of the outside of the tire, the tire has to be turned about its axis with respect to the lighting means and with respect to the cameras. The digital images obtained are then processed and compared against reference images in order to determine whether there might be any surface and appearance anomalies in the tire. For further details, reference may for example be made to patent applications EP-A2-1 959 227, EP-A1-2 023 078 and EP-A1-2 172 737.\nIn order to carry out such an inspection with a good degree of precision, use is generally made of an image acquisition device equipped with a laser lighting means that allows a line of light to be projected onto the tire and with a matrix camera capable of capturing the light reflected off the tire and which is oriented at a triangulation angle. The triangulation angle is the angle formed between the optical axis of the laser lighting means and the optical axis of the camera.\nIn order to be able to detect very small imperfections at the surface of the tire, the triangulation angle generally chosen is relatively large. As a result, the overall size of the image acquisition device is large.\nNow, in order to be able to capture a complete image of at least half the interior surface of the tire in a single rotation thereof, it is necessary to install a plurality of image acquisition devices penetrating the interior space of the tire.\nGiven the overall size of each image acquisition device it may therefore be particularly tricky if not to say impossible to achieve this installation for certain tire sizes, for example for tires in the current size range for passenger vehicles.\nThe present invention seeks to overcome this disadvantage."}
{"text": "The National Aeronautics and Space Administration (NASA) was developing new surface vehicles to support long range lunar exploration and the development of a lunar outpost. These vehicles will be heavier and will travel greater distances than the Lunar Roving Vehicle (LRV) developed for the Apollo program in the late 1960s. Consequently, new tires will be required to support up to ten times the weight and last for up to one hundred times the travel distance as compared to those used on the Apollo LRV, thereby requiring operational characteristics similar to passenger vehicles used on earth. However, conventional rubber pneumatic tires cannot function acceptably on the moon.\nFor example, rubber properties vary significantly between the cold temperatures experienced in shadow (down to 40 K) and the hot temperatures in sunlight (up to 400 K). Further, rubber degrades when exposed to direct solar radiation, without atmospheric protection. Finally, an air-filled tire is not permissible for manned lunar vehicles because of the possibility of a flat tire. To overcome these limitations, a tire design was developed for the Apollo LRV and was successfully used on Apollo missions 15, 16, and 17. This tire was woven from music wire, which was robust to lunar temperature variations and solar radiation, operated in vacuum, and did not require air for load support. This structure further functioned to contour to the lunar terrain, which facilitated traction and reduced vibration transfer to the Apollo LRV.\nBecause of the new weight and distance requirements for lunar vehicles, a tire with greater strength and durability would be desirable.\nOne conventional wheel and non-pneumatic tire assembly has a variable diameter which, in addition to changing its diameter, may also change its width, thereby increasing the area of the wheel that engages the ground. Thus, this non-pneumatic tire may be adjusted to increase a vehicle's performance according to the terrain over which it is traveling. This tire has arching members with first and second ends connecting a wheel hub. The arching members extend outwardly in an arc between the first and second ends. The arching members form a plurality of flexible hoops spaced circumferentially around the hub and extending radially outward from the hub.\nMore specifically, the conventional non-pneumatic tire forms a cage composed of thirty-eight equally spaced radially extending hoops that arch between axially outer rims of a hub. The hoops are made of helical steel springs filled by wires cut to a desired length and threaded through the center of the springs. The conventional hub may be expanded/contracted axially for varying the diameter of the tire.\nThe original wire mesh design of the Apollo LRV tire was found to not be readily scaleable. Specifically, the increase in wire diameter to create a tire that supported ten times the load of the original design created two significant limitations: 1) the ability to contour to the terrain was lost, thus limiting traction and ability to isolate vibration; and 2) the increased wire stresses limited functional life. A tire in accordance with the present invention overcomes these limitations, making the tire an innovative technological advance for Moon, Earth, and other planetary surfaces."}
{"text": "FIG. 1 illustrates an example of a conventional encoder for encoding multi-object or multi-channel audio signals. Referring to the drawing, a Spatial Audio Coding (SAC) encoder 101 is presented as an example of a conventional multi-object or multi-channel audio signal encoder, and it extracts spatial cues, which are to be described later, from the input signals, i.e., multi-object or multi-channel audio signals and transmits the spatial cues, while down-mixing the audio signals and transmits them in the form of mono or stereo signals.\nSAC technology relates to a method of representing multi-object or multi-channel audio signals as down-mixed mono or stereo signals and spatial cue information, and transmitting and recovering them. The SAC technology can transmit high-quality multi-channel signals even at a low bit rate. The SAC technology focuses on analyzing multi-object or multi-channel audio signals according to each sub-band, and recovering original signals from the down-mixed signals based on the spatial cue information for each sub-band. Thus, the spatial cue information includes significant information needed for recovering the original signals in a decoding process, and the information becomes a major factor that determines the sound quality of the audio signals recovered in an SAC decoding device. Moving Picture Experts Group (MPEG) based on SAC technology is undergoing standardization in the name of MPEG Surround, and Channel Level Difference (CLD) is used as spatial cue.\nThe present invention is directed to an apparatus and method for controlling rendering of multi-object or multi-channel audio signals based on spatial cue transmitted from an encoder, while the multi-object or multi-channel audio signals are down-mixed and transmitted from the encoder and decoded.\nConventionally, a graphic equalizer equipped with a frequency analyzer was usually utilized to recover mono or stereo audio signals. The multi-object or multi-channel audio signals can be positioned diversely in a space. However, the positions of audio signals generated from the multi-object or multi-channel audio signals are recognized and recovered uniquely to a decoding device in the current technology."}
{"text": "The invention relates to the technical field of power sockets, in particular to an anti-electric shock power socket.\nThe power sockets available on the market have multiple pressure units; each pressure unit usually includes a negative electrode holder and a positive electrode holder; the negative and positive electrode holders are connected to the power line by welding and wiring. In the viewpoint of manufacturing, the power socket is usually full of wires connected to each pressure unit, which not only makes the manufacturing inconvenient, it also increases cost. Besides, the power socket is internally wound with several segments of wires connected to each pressure unit, so overheating and short-circuiting phenomena are easily generated due to the inductive effect to cause danger. Thus, the existing power socket is necessary to improve because of the mentioned defects."}
{"text": "The present invention relates to an apparatus for taking care of the coat or fur of animals, and including a housing having an approximately rectangular section, a connector acting as a handle for connection to a suction device, a suction chamber and a comb.\nFew apparatus exist in the art for effecting treatments on the coat or fur of animals. Combs are known having some hollow teeth whereby it is possible to blow treating products onto the coat. Driers are also known which inject lukewarm air into the coat and are combined with a hair-gathering device. For the hair of human beings, French Pat. No. 1,536,475 discloses a comb associated with a suction case.\nThe major drawback of apparatus for treating the coat of animals resides in the fact that the chamber in which the treatment is effected becomes charged with dust, hairs or treating products with the resulting risk of allergy of the respiratory tracts of the animals. Further, if it is desired to effect a complete treatment, several successive operations must be carried out."}
{"text": "Lithographic printing typically involves the use of a so-called printing master such as a printing plate which is mounted on a cylinder of a rotary printing press. The master carries a lithographic image on its surface and a print is obtained by applying ink to said image and then transferring the ink from the master onto a receiver material, which is typically paper. In conventional lithographic printing, ink as well as an aqueous fountain solution (also called dampening liquid) are supplied to the lithographic image which consists of oleophilic (or hydrophobic, i.e. ink-accepting, water-repelling) areas as well as hydrophilic (or oleophobic, i.e. water-accepting, ink-repelling) areas. In so-called driographic printing, the lithographic image consists of ink-accepting and ink-abhesive (ink-repelling) areas and during driographic printing, only ink is supplied to the master.\nPrinting masters are generally obtained by the so-called computer-to-film method wherein various pre-press steps such as typeface selection, scanning, color separation, screening, trapping, layout and imposition are accomplished digitally and each color selection is transferred to graphic arts film using an image-setter. After processing, the film can be used as a mask for the exposure of an imaging material called plate precursor and after plate processing, a printing plate is obtained which can be used as a master.\nA typical photosensitive printing plate precursor for computer-to-film methods comprises a hydrophilic support and an image-recording layer which includes UV-sensitive compositions. Upon image-wise exposure of a negative-working plate, typically by means of a film mask in a UV contact frame, the exposed image areas become insoluble and the unexposed areas remain soluble in an aqueous alkaline developer. The plate is then processed with the developer to remove the diazonium salt or diazo resin in the unexposed areas. So the exposed areas define the image areas (printing areas) of the printing master, and such printing plate precursors are therefore called ‘negative-working’. Also positive-working materials, wherein the exposed areas define the non-printing areas, are known, e.g. plates having a novolac/naphtoquinone-diazide coating which dissolves in the developer only at exposed areas.\nIn addition to the above- photosensitive materials, also heat-sensitive printing plate precursors have become very popular. Such thermal materials offer the advantage of daylight-stability and are especially used in the so-called computer-to-plate method wherein the plate precursor is directly exposed, i.e. without the use of a film mask. The material is exposed to heat or to infrared light and the generated heat triggers a (physico-)chemical process, such as ablation, polymerization, insolubilization by cross-linking of a polymer or by particle coagulation of a thermoplastic polymer latex, and solubilization by the destruction of intermolecular interactions.\nThermal plates which require no processing are also known; such plates are typically of the so-called ablative type, i.e. the differentiation between hydrophilic and oleophilic areas is produced by heat-induced ablation of one or more layers of the coating, so that at exposed areas a surface is revealed which has a different affinity towards ink or fountain than the surface of the unexposed coating. A major problem associated with ablative plates, however, is the generation of ablation debris which may contaminate the electronics and optics of the exposure device and which needs to be removed from the plate by wiping it with a cleaning solvent, so that ablative plates are often not truly processless. Ablation debris which is deposited onto the plate's surface may also interfere during the printing process.\nOther thermal plates that require no processing are described in U.S. Pat. No. 5,855,173, U.S. Pat. Nos. 5,839,369 and 5,839,370 where a method relying on the image-wise hydrophilic-hydrophobic transition of a ceramic such as a zirconia ceramic and the subsequent reverse transition in an image erasure step. This image-wise transition is obtained by exposure to infrared laser irradiation at a wavelength of 1064 nm at high power (the average power is 1 W to 50 W and the peak power lies between 6 kW and 100 kW) which induces local ablation and formation of substoichiometric zirconia. U.S. Pat. No. 5,893,328, U.S. Pat. No. 5,836,248 and U.S. Pat. No. 5,836,249 disclose a printing material comprising a composite of zirconia alloy and α-alumina which can be imaged using similar exposure means to cause localized “melting” of the alloy in the exposed areas and thereby creating hydrophobic/oleophilic surfaces. A similar printing material containing an alloy of zirconium oxide and Yttrium oxide is described in U.S. Pat. No. 5,870,956. The high laser power output required in these prior art methods implies the use of expensive exposure devices.\nAnother type of processless plates are printing plates based on a so-called “switching” reaction where a hydrophilic surface is irreversibly changed into an oleophilic surface or vice versa by imagewise exposure. EP 652 483 for example, describes a positive working printing plate based on an acid catalyzed cleavage of acid-labile groups pendant from a polymer backbone. EP 200 488 and U.S. Pat. No. 4,081,572 describe negative working plates where a hydrophilic/hydrophobic conversion is obtained by a chemical reaction upon imagewise exposure to heat. Other examples of processless plates are based on the thermally induced rupture of microcapsules and the subsequent reaction of the microencapsulated oleophilic materials (isocyanates) with functional (hydroxyl-)groups on cross-linked hydrophilic binders (U.S. Pat. No. 5,569,573; EP 646 476; WO94/2395; WO98/29258).\nU.S. Pat. No. 6,582,882 describes an imaging element comprising a graft copolymer having a hydrophobic backbone and a plurality of pendant hydrophilic groups or a plurality of pendant groups comprising hydrophilic and hydrophobic segments. Upon exposure of the imaging element to thermal energy, the exposed areas become less soluble in a developer than the unexposed areas.\nU.S. Pat. No. 6,362,274 describes grafted copolymers comprising three sequences: one sequence for anchoring on solid particles such as pigments and fillers, one hydrophobic sequence and one hydrophilic sequence for using the copolymers in aqueous and/or organic medium. The disclosed copolymers are of particular interest in a wide range of paint formulations; there is no reference in the cited prior art document to lithographic printing plates.\nNone of the prior art discloses the heat-sensitive copolymer of the present invention in lithographic printing plates."}
{"text": "The term “wave energy” is used herein to include radiation as well as wave energies transmitted by various mediums, and embraces electromagnetic waves or radiations; sonic, supersonic, and ultrasonic waves; neutrons, protons, deuterons, and other corpuscular radiations. The term “electromagnetic waves” includes, e.g., X-ray, gamma-ray, ultraviolet, infra red, and visible light rays, and short electric and radio waves. These definitions and terms are consistent with those used by the U.S. Patent and Trademark Office for classification purposes.\nElectromagnetic radiation (EMR) is one of the most pervasive forms of wave energy known and used by man. Sunlight, both within and beyond the visible spectrum, is one example of EMR wave energy that has been highly beneficial to man, and all life on earth. Man has harnessed the benefits of EMR through the development of such great inventions as the light bulb, X-rays and radio waves. Another utilization of wave energy, welding, has been an important factor in the development of modem technologies. Similarly, the steel industry grew rapidly due to another wave energy technology known as the electric arc, used for melting iron ore and converting it to steel.\nAlthough the industrial revolution has brought about the luxuries and simplicities of life there has been a world-wide downside—pollution. Nowadays, it is common to find many waterways, drinking water sources, air and soils contaminated with pollutants. Although many technologies have been developed for removing contaminants from water and air, more and more engineers and scientists are turning to ultraviolet (UV) light systems for treating fluids.\nOne of the greatest inventions of all time is the light bulb. One of the earliest forms, Thomas Edison's carbon arc light bulb, has all but faded into extinction. A few very large carbon arc systems such as World War II vintage carbon arc searchlights are available today as rebuilt units. However, in general carbon arc lights are not utilized today for industrial and residential lighting, as searchlights or for photochemical reaction type applications. The carbon arc light bulb has been almost entirely replaced for those purposes by fluorescent bulbs, high pressure mercury vapor lamps, compact ceramic lamps and high pressure xenon lamps.\nAnother form of wave energy is sonic energy, and particularly ultrasonic energy. Ultrasonic waves are more commonly referred to as ultrasound. Ultrasound plays an important role in medical diagnostics, submarine sonar, ship sonar, non-destructive testing of metal, and cleaning equipment.\nPhotochemical reactions are well known and well documented. The use of electromagnetic radiation (EMR) particularly within the ultraviolet (UV) region, between 4 to 400 nanometers (nm), for treatment of fluids such as disinfection of drinking water and wastewater, free radical generation (hydroxyl radicals, chlorine radicals, etc.) and removal of noxious air contaminants such as VOCs, NOx, SOx from flue gas (off gas, tail gas) has gained in popularity over the past decade. In addition, advanced UV processes such as Advanced Oxidation, which incorporates UV light with an oxidant such as ozone or hydrogen peroxide, generate free radicals for decomposition of contaminants. Another UV art, which is rapidly receiving attention, is PhotoCatalytic Oxidation (PCO). PCO technology incorporates a semi-conductor catalyst in combination with UV radiation for generation of free radicals. PCO technology can use sunlight since the photons of interest for the most common photocatalyst, TiO2, lies between 320 to 400 nm and more specifically at about 365 nm.\nAlthough there are many methods and devices known in the art, their applications are specific and limited. It is highly unlikely, for example, that an EMR device used for water disinfection can be used for air purification. Likewise, a PCO device designed for air disinfection has significant drawbacks when an attempt is made to utilize that technology in a liquid environment. This can be clearly demonstrated by a close review of the prior art.\nThere are a number of major obstacles which have not been overcome by the prior art. One major obstacle is residence time versus absorbance. The result of the failure of the prior art to overcome this obstacle is perceived by the end-user (customer) as a problem of “not having enough lights,” of “light penetration,” and/or of “excessive energy requirements (inefficient)”. When a photochemical reactor is designed and built for a specific set of parameters, the actual parameters in which the reactor is operated are variable and often change after the reactor put into service. Operators must make do with the photochemical reactor or the reactor is taken out of operation and decommissioned.\nGiven the past inadequacies of wave energy systems, in particular EMR systems employing UV/visible radiation, many designs have incorporated more lights, transparent glass sleeve wipers, redundant systems, or methods, which allows for more/better contact between the contaminants and the photons or the photocatalyst and the photons. However, the changes still do not overcome the problem of residence time versus absorbance. Thus, if the problems associated with both residence time and absorbance can be eliminated or solved, then an ideal photochemical reactor as well as method (for photochemical reactions) can be designed for many different applications. The problem of absorbance is inversely proportional to EMR transmission.\nA number of attempts to advance the technology and to overcome problems and drawbacks in this field have been made, and are reflected in a number of patent documents.\nU.S. Pat. No. 3,998,477—dated Dec. 21, 1976, discloses a device for non-rigid connection of two rigid cylindrical pipes, which comprise the combined use of flexible double lipped gaskets and toroid flexible gaskets surrounding a bulb-shaped zone formed in one of the pipes. The device is particularly useful for fastening fragile tubes containing light emitters to metal reactors used for photochemical processes.\nU.S. Pat. No. 4,002,918—dated Jan. 11, 1977, discloses an apparatus for the irradiation of fluids in which the fluid is conducted along the walls of a container having walls which are permeable for the radiation to which the fluid is exposed. Radiation sources are arranged around the container and an active rotor is disposed within the container. The rotor consists of a body having axial bores and pins movably disposed in the bores and adapted to engage with their front ends the container walls thereby to wipe any deposits from the container walls during rotation of the rotor.\nU.S. Pat. No. 4,317,041—dated Feb. 23, 1982, discloses various embodiments of photo-reactors in which there are at least two radiation chambers with a window arranged therebetween. UV radiation is introduced into one of the chambers at a side opposite the window so that it passes through that chamber, through the window and into the other chamber. The fluid medium to be purified is passed through the chambers and subjected to the radiation while in the chambers. The flow of the medium is through the chambers in series in some embodiments and in parallel in others. An embodiment is disclosed wherein a recirculation line is established around the reactor with the recirculation being continuous or intermittent. When intermittent the purified fluid medium also is drawn off intermittently, between the periods of recirculation. In some embodiments, the amount of radiation traversing all the chambers is monitored. If the monitored amount drops below a given amount, the apparatus is shut down. Alternatively, the rate of flow of the medium is adjusted, based on that monitored amount, with the rate of flow increasing or decreasing, respectively, in response to increases or decreases in that amount.\nU.S. Pat. No. 4,476,105—dated Oct. 9, 1984, relates to a process for producing gaseous hydrogen and oxygen from water. The process is conducted in a photolytic reactor which contains a water-suspension of a photoactive material containing a hydrogen-liberating catalyst. The reactor also includes a column for receiving gaseous hydrogen and oxygen evolved from the liquid phase. To avoid oxygen-inactivation of the catalyst, the reactor is evacuated continuously by an external pump which circulates the evolved gases through means for selectively recovering hydrogen therefrom. The pump also cools the reactor by evaporating water from the liquid phase. Preferably, product recovery is effected by selectively diffusing the hydrogen through a heated semipermeable membrane, while maintaining across the membrane a magnetic field gradient which biases the oxygen away from the heated membrane. This promotes separation, minimizes the back-reaction of hydrogen and oxygen, and protects the membrane.\nU.S. Pat. No. 5,126,111—dated Jun. 30, 1992, discloses a method of removing, reducing or detoxifying organic pollutants from a fluid, water or air, by contacting the fluid with a photoreactive metal semiconductor material in the presence of ultraviolet light of a wavelength to activate the photoreactive material. This is improved by simultaneously contacting the photoreactive material with a substance that accepts electrons and thus inhibits hole-electron recombination. Such substance will be such as to readily accept electrons either from the conduction band or from superoxide ions, and to rapidly dissociate into harmless products.\nStill other photoreactors are described in U.S. Pat. Nos. 3,567,921; 3,769,517; 3,924,246; 4,296,066; 4,381,978; 4,454,835; 4,488,935; 4,544,470; 4,774,026; 4,863,608; 4,868,127; 4,957,773; 5,045,288; 5,094,815; and 5,149,377.\nU.S. Pat. No. 5,439,652 (Sczechowski, et al.) issued on Aug. 8, 1995 states that a Beer's law type expression was found for the transmitted light as a function of the TiO2 loading. From this relationship, the calculated light penetration depth for the 0.4% (by weight) Degussa TiO2 slurry used in these experiments was approximately 1 mm.\nU.S. Pat. No. 5,994,705 (Cooke, et al) issued on Nov. 30, 1999 is a continuation of U.S. application Ser. No. 08/946,647, filed on Oct. 7, 1997 now U.S. Pat. No. 5,866,910 which is a continuation of U.S. application Ser. No. 08/438,234, filed on May 9, 1995, now U.S. Pat. No. 5,696,380 discloses a flow-through photochemical reactor includes a reactor body, which circumscribes a longitudinally extending channel having a generally annular cross section. This channel accommodates fluids passing between an inner wall of the reactor body and an outer wall of a photon-transmitting tube that is housed internally thereof. In addition, the reactor includes mechanically static, fluid-dynamic elements for passively inducing substantial turbulent flow within a fluid as it passes through the channel. This arrangement substantially increases the uniformity of the fluid's exposure to photons radiating from a source within the tube into the fluid and it is conducted through the channel.\nCalgon Corporation's U.S. Pat. No. 6,565,803 issued to Bolton, et al. on May 20, 2003 and titled, “Method for the inactivation of cryptosporidium parvum using ultraviolet light,” has a major drawback. The system utilizes mercury vapor bulbs housed in a quartz tube. Mercury is a pollutant that is transferred via the food chain. Any UV system incorporating a “bulb” is prone to burn out. Further, the glass or quartz envelope and the bulb become solarized due to the UV light. Consequently, dosages as specified within the '803 patent may not be sufficient to inactivate cryptosporidium.\nMany other types of wave energy apparatuses are known in the prior art, but none of the known prior art utilizes the approaches encompassed within the scope of the present invention."}
{"text": "1. Field of the Invention\nThe invention relates to clutch units of the oil shear type, and particularly to units having interleaved clutch plates and discs which are alternately clamped together or separated in an axial direction. The invention is directed toward a very large size clutch unit such as would be found in an aircraft catapult system for selectively winding or tensioning a nylon belt or tape which is in turn connected to a carriage secured to the aircraft to be launched or catapulted.\n2. Description of the Prior Art\nMy U.S. Pat. No. 3,696,898 issued Oct. 10, 1972 shows a clutch-brake unit in which a plurality of interleaved clutch discs and plates are provided, all of which are simultaneously actuated by a single piston. The arrangement shown in this patent would be unsatisfactory for the purpose of actuating large size clutch units which must transmit high horsepowers. Among the problems to be dealt with in the construction of large size oil shear type clutch units are the need for maximum application of forces to couple the plates and discs, the cooling requirements for the part, and most importantly, the need for minimizing residual drag. To illustrate the magnitude of residual drag problems, it should be noted that in an 85,000 horsepower clutch operating at 1200 R.P.M., the residual drag is in the order of 8-10,000 horsepower. The residual drag in the apparatus of the present invention, on the other hand, is less than 500 horsepower."}
{"text": "Advertising and marketing displays often consist of non-static images, such as video images and other digital displays. For example, large-scale digital displays are widely used in major urban centers, as well as in stadiums, arenas, exterior walls of buildings and other applications. There exists an increasing need for smaller scale displays that are capable of displaying a static or dynamic video image so as to better attract the attention of onlookers. Displays of this nature may be used for advertising, public service information, and the like. As used in this patent specification, the terms “video display” or “flat panel display” encompasses any type of image display, whether for commercial advertising or other purposes. Preferably, such a display is a digital static or dynamic video display, but other video display types are known.\nThere is a need to provide video display systems that are suitable for use as a street or floor-level display for use in a harsh environment, such as outdoor use in locations where the units may be exposed to direct sun, and other situations of extreme heat or cold. Flat panel display screens, such as LCD panels, provide a high degree of resolution and are well-suited for displaying advertising and other images on a smaller scale, for example street-level displays. As well, advances in LCD and other flat panel screen technologies have permitted very bright display capabilities, making this technology in principal suitable for outdoor use in direct or near-direct sunlight. However, apart from certain closely monitored environments such as airports and retail stores, exposed video units can be exposed to vandalism and theft. As well, outdoor use of an exposed screen is usually difficult as a result of exposure to moisture and extreme temperature conditions. One solution is to enclose a video screen within a housing or enclosure. An enclosure may also find uses in indoor settings, where for security or other reasons it is desired to encase the video display unit within a robust enclosure.\nConventional LCD monitors are suitable for operation in a temperature range of approximately 5° C. to 40° C. At temperatures above this range, the screen may not function and may be temporarily or permanently damaged, while below this range the display may also not function properly. This presents a drawback for outdoor usage of LCD displays (and other flat panel digital monitors) that are enclosed within a housing, since temperatures can exceed this range. The interior of such an enclosure can experience a greatly elevated temperature, resulting from the combination of internally-generated heat from the video monitor, as well as absorbing of heat from the sun. There is therefore a need to provide a system which includes a secure enclosure for a flat panel digital display unit, but which is capable of maintaining a suitable temperature range for a video unit. Thermal management of the system should be provided in a manner which is efficient, reasonably quiet and reliable within a range of conditions."}
{"text": "As electronic mail becomes more popular, users face large volumes of messages in their electronic mail folders. This makes at least two important tasks difficult: triage of new messages and searching for existing messages. Whether for searching or triage, a user would like to find or dispose of the messages quickly, but may not know an optimal strategy.\nFor example, a user with 100 new electronic messages may scan the messages one by one, spending too much time addressing less relevant messages, and taking longer to discover the more important messages. Similarly, if a search query returns many messages, the user may need to scan the messages individually to find the desired one(s). In either case, the user spends too much time looking at less relevant messages."}
{"text": "Liquid crystals are applied to display media in which the reversible movement of liquid crystal molecules is made use of, such as display devices of various modes typified by twisted nematic (TN), vertical alignment (VA), and in-place-switching (IPS). Besides the application to display media, the liquid crystals, particularly those having a polymerizable functional group have been studied for applicability to optically anisotropic materials, such as a retardation film, a polarizer, a polarizing prism, a luminance-improving film, a low pass filter, various optical filters, and optical fibers, taking advantage of their anisotropy in physical properties, such as refractive index, dielectric constant, magnetic susceptibility, elastic modulus, and thermal expansion coefficient, as well as their alignment properties. It is important for the optically anisotropic material (polymer) obtained by the polymerization to have not only the optical anisotropy but other characteristics, such as polymerization rate, transparency, mechanical strength, coating properties, solubility, crystallinity, shrinking properties, water permeability, water absorption, melting point, glass transition point, clear point, chemical resistance, and heat resistance.\nThe optically anisotropic material (polymer) is obtained by, for example, uniformly aligning the molecules of a liquid crystal compound having a polymerizable functional group or a polymerizable composition containing the liquid crystal compound into a liquid crystal phase and irradiating the compound or the composition being in the liquid crystal phase with energy rays, such as ultraviolet rays, to cause photopolymerization. It is required to fix the aligned state of the liquid crystal compound uniformly and semi-permanently.\nWhen the polymerizable composition has a high liquid crystal phase transition temperature, photopolymerization induced by energy rays may unintentionally invite thermal polymerization, which disturbs the uniform alignment of the liquid crystal molecules, making it difficult to fix a desired state of alignment. In order to facilitate temperature control during cure, a polymerizable composition showing a liquid crystal phase at or near room temperature is demanded.\nThe polymer is obtained by polymerizing the polymerizable composition in the form of coating film applied to a substrate. If the composition contains a non-polymerizable compound, the resulting polymer film may have insufficient strength or contain residual stress-induced strain. Removing a non-polymerizable compound using, e.g., a solvent can result in a failure to retain film homogeneity and cause unevenness. To obtain a polymer film with a uniform thickness, it is therefore preferred to apply a polymerizable composition in the form of a solution in a solvent to a substrate. Hence, it is desirable for a liquid crystal compound or a polymerizable composition containing it to have good solubility in a solvent.\nTo cope with the tendencies to reduction in thickness and weight of display devices, components making up the display devices including optical films have been requested to have a reduced thickness.\nA polymerizable liquid crystal compound having an acryloyloxy group as a polymerizable functional group exhibits high polymerization reactivity, i.e., readily polymerizes through a convenient method (e.g., UV irradiation), and the resulting polymer has high transparency and has therefore been extensively studied for use as an optically anisotropic material (see, e.g., Patent Document 1, Patent Document 2, Patent Document 3, and Patent Document 4). However, any of the polymerizable compositions disclosed in these reference documents has the disadvantage of difficulty in controlling film formation with a reduced thickness.\nIn general, formation of a thick film from a polymerizable composition containing a liquid crystal compound encounters with difficulty in controlling molecular alignment of the liquid crystal compound, which can cause problems, such as reduction in transmittance and coloration. On the other hand, a thin polymer film having satisfactory molecular alignment over the entire area can be obtained, but formation of a thin film has difficulty in thickness control, readily resulting in non-uniform surface condition or crystallization. Furthermore, the liquid crystal state resulting from alignment control shows poor stability and tends to be disturbed before curing by irradiation with, e.g., UV light. Conventional known polymerizable compositions have thus failed to provide a satisfactory polymer.\nAlthough Patent Document 5 supra discloses an epoxy compound that cures on irradiation with light, it is not pertinent to utilization as a liquid crystal material, nor does it teach or even suggest the effect on control in film formation and fixation of liquid crystal alignment.    Patent Document 1: JP 2006-193596A    Patent Document 2: JP 2005-196221A    Patent Document 3: JP 11-148080A    Patent Document 4: JP 8-231958A    Patent Document 5: JP 2006-84985A"}
{"text": "The present invention relates generally to a patient support and, more particularly, to a patient sensor for detecting the presence of a patient supported on a patient support surface.\nAccording to an illustrative embodiment of the present disclosure, a patient sensor is configured to detect a force generated by a patient supported on a patient support surface, the patient sensor comprising a base, a detector operably coupled to the base and positioned below the patient support surface, wherein the detector is configured to detect a force applied to the patient support surface. A force collector extends outwardly from the base toward the detector and is configured to collect and concentrate the force applied to the patient support surface and to direct the force to the detector.\nAccording to a further illustrative embodiment of the present disclosure, a patient sensor comprises a base, a force sensing resistor having a resistance value configured to change depending upon the amount of force applied thereto, and a shear reducing coupler operably coupling the base to the force sensing resistor such that at least portions of the force sensor resistor are free to move relative to the base.\nAccording to yet another illustrative embodiment of the present disclosure, a patient sensor is configured to detect force exerted by a patient against a patient support surface, the patient sensor comprising a base, and a force sensing resistor operably coupled to the base, wherein the force sensing resistor has a resistance value configured to change depending upon the amount of force applied thereto. A spacer is supported by the base and is configured to position the base in spaced relation to the force sensing resistor.\nFurther illustratively, the force sensing resistor is operably coupled to a controller and includes a fault condition configured to be detected by the controller. Illustratively, a resistor is coupled intermediate the force sensing resistor and the controller for providing the fault condition.\nAdditional features and advantages of the present invention will become apparent to those skilled in the art upon consideration of the following detailed description of illustrated embodiments exemplifying the best mode of carrying out the invention as presently perceived."}
{"text": "The present invention concerns a hot air circulating oven range using shape memory alloy springs to drive the damper which changes over the airflow inside the oven.\nShape memory alloys try to return to their original shape when they reach a certain temperature, which enables the use of this material as a spring to drive the opening/closing of a damper depending on temperature.\nHowever, with the normal type of shape memory alloy, the force with which the alloy trys to return to its original shape reduces the effective usage or length of use of the alloy, or in other words the alloy deteriorates, so that when a bias spring is used to keep a damper open, the force with which the shape alloy must pull against this spring to close the damper gradually reduces as it deteriorates, and eventually the force of the shape memory alloy is so small that it cannot overcome the force of the bias spring for closing the damper. This in turn results in a lower thermal efficiency when cooking with the hot air circulation method."}
{"text": "A wire cable tray generally has an elongate form and is of U-section. It comprises longitudinal wires, also termed warp wires, and transverse wires, also termed weft wires. The warp wires are generally straight and the weft wires are U-shaped. The weft wires are welded to the warp wires and are regularly spaced. The cable tray so produced has a bottom adapted to serve as a support for the cables or the like, and side walls, or flanges, adapted to hold the cables on the bottom.\nIn all that follows, for the sake of simplicity and clarity, the description will be made with reference to electric cables. However, the invention applies to any other type of longilinear item that can be supported by a wire cable tray: tube for conveyance of fluid (pneumatic or hydraulic), optic fiber, etc.\nSuch cable trays are appropriate anywhere that it is desired to pass electric cables. However, when a cable tray is adapted to support electric cables at the floor, it is preferable to raise it in order to protect the electric cables which it conveys, in particular in case of flooding.\nTo raise a cable tray, it is known for example to have feet regularly distributed along the cable tray. For example, there may be two feet on the same weft wire and this may be repeated at regular intervals. It is also known to rest the cable tray on a stirrup fixed to the floor. The connection between the stirrup and the cable tray is made for example using hooks formed on the stirrup which come into engagement with a warp wire of the cable tray.\nTo splice together two cable trays, it is known to use for example a splicing member for linking two cable trays placed end to end. Two weft wires are then disposed face to face. The splicing member comprises for example clipping means enabling it to clip onto each of the weft wires that face each other, so providing the connection between the cable trays.\nThe devices known from the prior art are not intended to provide at the same time the raising and the splicing together of two cable trays. Furthermore, the known raising accessories most often cannot be mounted in advance on the cable tray for which they are intended. When such a pre-mounting is possible, this is carried out using nuts and bolts or welds. Such a pre-mounting is costly in terms of time."}
{"text": "The present invention disclosed herein relates to a signal information transmitting device and a signal information transmitting method.\nA car having an Electronic Control Unit (ECU) is driven according to a control of the ECU. Exemplarily, the car may control a fuel injection amount, front wheel and rear wheel brakes, and front wheel and rear wheel driving powers according to a control of the ECU. Additionally, the car may control whether to drive the engine according to a control of the ECU. For example, when the car waits for a traffic light, the engine of the car may be stopped according to a control of the ECU. Also, when the car starts, the engine of the car may be driven according to a control of the ECU.\nAccording to the timings that the ECU stops and drives the engine, the fuel consumption efficiency and operating performance of the car may be changed. For example, as a time until the engine stops after the car stops is longer, the fuel consumption efficiency of the car is decreased. Moreover, as a time until the engine starts after a start signal is turned on is shorter, the operating performance of the car is increased. That is, when the ECU of the car is synchronized to a signal switching timing of a traffic light, the fuel consumption efficiency and operating performance of the car may be optimized."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nRemote access appliances, also known as “Keyboard/Video/Mouse” (“KVM”) appliances, are often used to enable remote access from a user's personal electronic device, typically a laptop or a desktop computer, to a target device within a data center. Often the target device is a server, but this is only one example, and the target device could be any data center device that the user needs to communicate with.\nTypically the KVM appliance is coupled to the target device, which as noted above is often a data center device, and more typically a server. The coupling is achieved by an interface device known in the industry as a “rack interface pod” (“RIP”) or a “digital rack interface pod” (“DRIP”). In either case, the interface device typically connects to an RJ45 port of the remote access appliance via a first cable, which may be a CAT5 or similar cable. A second cable, which is typically a video cable which is hard-wired to the interface device, extends out from the interface device and is coupled to a video port of the target device (e.g., a server). A third cable, typically a USB cable, may be hardwired or coupled via a releasable USB connector to the interface device. The third cable forms a serial communications link between the interface device and the target device. Thus, the second and third cables form the means by which serial communications and video information may be transferred between the interface device and the target device.\nWhile the above described interface device has performed well in a data center environment, the requirement of three cables extending from the interface device (i.e., CAT5, USB and video) can make for a complex arrangement of cabling at the equipment rack where the interface device is being used. If an entire equipment rack is filled with servers each having an interface device coupled thereto, one can easily see that this creates a voluminous arrangement of cables that must be managed. If all of the equipment racks in a row are filled with servers or other devices to which such interface devices are coupled, the amount of cabling grows almost exponentially. Accordingly, a need still exists for an interface device that is able to interface data center devices with a KVM appliance with a significantly reduced amount of cabling."}
{"text": "Organizations use storage virtualization systems to manage data objects stored on multiple network attached storage (“NAS”) systems. A storage virtualization system can migrate data objects from one NAS system (the “source,” “source server” or “src”) to another NAS system (the “destination,” “destination server” or “dst”). A storage virtualization system can also synchronize data objects between two NAS systems. In order to perform these operations, a storage virtualization system will typically be placed as an intermediary between the NAS system that stores the data objects and the client computer that accesses or requests access to the stored data objects. The client computer, storage virtualization system and the associated NAS system may be connected over a network. The storage virtualization system may begin performing its migration or synchronization operations once a second NAS system is added to the network. Storage virtualization provides a mechanism for real-time or near real-time data object management without requiring a shutdown or disconnection of NAS resources, or any intervention from users. As such, storage virtualization provides an effective and efficient way to manage data objects stored on multiple NAS systems without causing downtime or reducing productivity.\nIn order to efficiently perform storage virtualization operations, the storage virtualization system must be able to identify and locate the data objects stored on each of its associated NAS systems. For example, the storage virtualization system must monitor all of the data files, data directories, hard links, soft links and devices on its associated NAS systems. The storage virtualization system may therefore internally store a data structure containing information or “metadata” about the data objects stored on its associated NAS systems. The data structure may be a mapping index, a database, a File Identifier (“FID”) table, or directory tree that identifies each data object and its location on the associated NAS systems. The data structure may also include relevant path address, file attributes, timestamps or other such information for those data objects. Metadata in the data structure may originate from the associated NAS system that stores the data object, or the storage virtualization system that manages the data object. Using such a data structure enables the storage virtualization system to have an updated “snapshot” of the data objects it is managing. Changes made to those data objects, such as renaming or editing data object content, may be reflected in the data structure.\nDue to hardware and performance constraints, a storage virtualization system may only have a finite amount of internal physical memory in which to store the data structure. For example, a storage virtualization system with 16 gigabytes of internal memory may only be able to store information on 50-60 million data objects. As a result, the storage virtualization system may only be able to manage 50-60 million data objects. NAS system manufacturers are continuously increasing the memory and data object storage capacity of their NAS systems. As the number of data objects that can be stored on these NAS systems grows, a storage virtualization system associated with these larger NAS systems will need to keep up with the increased storage demands.\nWhat is therefore needed is a way to increase the storage capacity of present and future storage virtualization systems so that they can manage larger numbers of data objects. The data objects themselves may be stored on a single NAS system or on multiple NAS systems associated with a storage virtualization system. What is further needed is a way to optimize the storage virtualization system's data object management capabilities without affecting performance or resulting in data loss."}
{"text": "1. Field of the Invention\nThe present invention relates to a learning-type classifying apparatus and a learning-type classifying method.\n2. Description of the Related Art\nAs a learning-type classifying apparatus, the configuration of an apparatus is shown in Jpn. Pat. Appln. KOKAI Publication No. 8-21803 wherein a neural network is so constructed that characteristic values of the defective portion extracted from within an image is input as an input pattern, and the kind of defect is output by learning, thereby classifying the defect.\nFurthermore, Jpn. Pat. Appln. KOKAI Publication No. 11-344450 shows a configuration which assists a teacher data creating operation, a configuration which uses statistical discrimination analysis to classify defective images on the basis of the created teacher data, etc. It is to be noted that the teacher data is a set of characteristic value information calculated from image information on or an image of a classification target and information on a correct class (=correct category) of the target (see FIG. 4 in the above-mentioned publication)."}
{"text": "Systems and methods have been used to analyze large amounts of text. Some schemes attempt to locate themes common to the text being analyzed and may thereafter associate respective text documents with appropriate themes. These systems may operate to separate information of interest from information that is less relevant or useful. Text analysis systems and methods are becoming increasingly useful as the volume of information available for analysis continues to increase. Some conventional arrangements suffer from drawbacks of being relatively inflexible to accommodate input of a user when performing analysis.\nAt least some aspects of the present disclosure provide systems, methods and articles of manufacture which improve user interactivity with respect to text processing. Additional aspects provide increased flexibility with respect to text being analyzed as discussed in further detail below."}
{"text": "1. Field of the Invention\nThis invention relates to a device enabling a means of maintaining a desired orientation of a means of catching a fish, such as, but not restricted to, a hook, snare, artificial bait, or lure.\n2. Background Information\nThroughout history, man has struggled to advance the art of catching fish, for various and sundry reasons such as, but not restricted to, hunger, sport, sand assurances that he is smarter than the fish.\nIn using jig hooks, and related hardware, existing art has problems in that there is no sure way to maintain a hook at a desired angle to a line prior to a fish taking the hook and also inside a mouth of a fish. The hooks in the prior art tend to either fail to reliably, snag the fish and/or they can become lodged deep in the mouth of the fish where removal of the hook is difficult and catch and release of the fish is also difficult.\nU.S. Pat. No. 2,592,664 to De Mello shows a prior art fishing lure. Primarily intended as a double hook, FIG. 1, '664 does show in FIG. 4 a single hook 29 with an upper eye 27 a lower eye 25, a lever arm 26 between them and a hook shaft 28 held perpendicular to a line 10, and a weight 11. The De Mello patent is primarily designed to maintain the hook in a weedless (column 2 line 40) orientation and is designed to snag fish in the lower jaw. Since most fish dive down when taking bait and hook, and since the fisherman will pull up to set the hook, trying to hook the lower jaw (column 2, line 38) is the least desirable target in many cases. The De Mello patent contemplates only a single position with the hook shaft 28 perpendicular to the line 10 (column 2 line 30) and the hook end 31 in a downward orientation. Finally, the hook 29 of DeMello has a very long hook shaft 28, the hook shaft 28 being 4 or 5 times longer than the apparent diameter of the hook curve 30 of the hook 29. What the long hook shaft 28 means is that the hook end 31 is likely to hook the fish deep in the mouth or gut making catch and release impossible.\nAs will be seen, the present invention overcomes the limitations of the prior art."}
{"text": "For many years paper greeting cards have been widely used for celebratory occasions such as birthdays, graduations, weddings and other commercial purposes. Traditional text information is generally found on paper greeting cards. More recently, sound has been added to traditional paper greeting cards to increase the personalization of the cards by delivering an audio message that is electronically embodied in circuitry that is carried within the greeting card. The ability to use sound in combination with printed matter such as with conventional printed greeting cards significantly enhances the communicative value of social and relational greetings. The availability of small voice recording sound modules has made sound-generating greeting cards increasingly popular. Small lighting systems have also been incorporated into social greeting products and novelties, and combined in circuits with sound and other features."}
{"text": "Account cards are used in many aspects of daily life, including gaining access to secure locations, providing identifying information, and performing monetary transactions. In some situations, mobile devices such as smartphones are replacing traditional cards, such as mobile devices having apps for making payments at point-of-sale locations or accessing a hotel room. But traditional cards still remain the most commonly-used mechanism for providing account information, thanks to their small form factor, convenience, and the highly developed infrastructure for receiving account information from cards such as magnetic card readers.\nCompared to mobile electronic devices, traditional cards have more inherent security risks. While mobile electronic devices are usually locked with a passcode or PIN to prevent unauthorized usage, traditional cards are inherently unlocked and can be used by anyone in possession of the card. Thus, when a card is lost or stolen, the finder or thief can use the card at will, until an intervening event occurs, such as the card issuer disabling the card, a diligent cashier refusing the card upon checking the cardholder's identification, or destruction of the card. Often times, card issuers and the true cardholders suffer the consequences of fraudulent activities by the time the card number is cancelled or changed.\nIn view of the shortcomings of current systems and methods for providing account information, a robust, secure, and efficient mechanism for controlling activation of account card devices is desired."}
{"text": "The present invention relates to a device for auto-adaptative direction and polarization filtering of radio waves received on a network of aerials coupled to a receiver.\nThe invention is applied in particular to the construction of an interference elimination device for receivers of electromagnetic waves travelling through the HF ionospheric channel.\nIt is known that a tactical high frequency reception aerial must in particular enable any user to establish radio connections with various transmitting stations, and to simultaneously hear the transmissions from the other users of the high frequency channel.\nBearing in mind the congestion of the high frequency channel, the receiver must be able to reject the signals coming from sources of interference, be they intentional or not.\nHowever, the known tactical high frequency adaptative aerial networks, generally formed from several suitably spaced vertical whips, do not enable exploitation of the polarization of the ionospheric waves. Whilst the performance levels expected from an adaptative aerial formed from three orthogonal dipoles, which does not have the shortcomings of the whips, are described in an article by RTE COMPTON entitled \"The tripole antenna: an adaptative array with a full polarization flexibility\" (IEEE trans antenne and propagation. Vol.: AP 29 No. 6 Nov. 81), the aerial which is described therein has the tactical disadvantage of having to be placed at the top of a mast so that the radiation patterns of the horizontal dipoles are correct. Moreover, as it is technologically very difficult to suppress the high frequency currents circulating in the external parts of the coaxial cables connecting the dipoles of the aerial to the receiver, the performance levels expected of this aerial configuration are, in practice, not obtained."}
{"text": "Snowmobiles are vehicles designed primarily to move or travel over snow covered terrain. Generally, snowmobiles are open air vehicles with a motorcycle style seat that can accommodate up to four people. Snowmobiles are usually powered by an internal combustion engine that drives a long continuous or endless track which is located in a rearward position of the snowmobile, and substantially below or underneath the seat. This continuous or endless track, powered by the snowmobile's motor, is configured to engage the snow covered terrain and propel the snowmobile over the snow. The track provides a similar function for the snowmobile as a tire does on a road vehicle.\nSnowmobiles are typically further equipped with a pair of skis positioned in a forward position. The function of the skis is to provide a front ground contact for stabilization, as well as to facilitate guiding or steering of the snowmobile. A pair of handlebars, accessible by an occupant on the seat, can be pivoted in order to turn the skis and thus turn or steer the snowmobile.\nThe front skis on a snowmobile, however, have several disadvantages that affect snowmobile performance and operation. For example, since the skis are designed to slide over the snow, they provide very little traction to the snowmobile. Consequently, the rear track often partially overpowers the traction of the ski and continues to push the snowmobile forward instead of through a turn. This partial overpowering of the ski traction causes the vehicle to make a wider turn, thereby affecting the cornering ability and overall handling of the snowmobile. Sometimes the rear track completely overpowers the skis and moves the snowmobile in a forward direction regardless of the direction the skis are turned. This is particularly the case when trying to traverse across hard surfaces such as hard pack snow, asphalt, or the like. The problem is compounded as snow provides a poor traction surface, and in the case of unpacked snow, is easily pushed or displaced. As such, when turning or moving sideways along a hill, the rider must often lean dramatically on the snowmobile in order to weight the front end and get the skis to bite into the snow. To account for these problems, those skilled in the art have equipped snowmobile skis with rails that extend down from the underside or contact portion of the skis. Although the rails are designed to increase the traction of the snowmobile by digging into the snow, these rails result in only marginal improvements.\nAnother disadvantage of the front skis is that they provide no assistance to the snowmobile's ability to accelerate, slow down, or reverse. In other words, they are passive and provide no driving force. In fact, since all the skis do is slide on or across the snow, they actually decrease the efficiency of the snowmobile as they tend to produce drag and a resistance that the rear endless track must overcome. Moreover, since the rear endless track of typical snowmobiles can sometimes become bogged down by snow, ice, rocks, and other debris becoming lodged in the track tunnel between the track and the seat, the additional resistance of the skis can stop the snowmobile from moving at all.\nStill another disadvantage of the front skis is that they tend to ride on the top or upper surface of the snow, whereas the rear track tends to dig into the snow due to its driving force as powered by the snowmobile's motor, thus actively displacing the snow. The result of this is that the snowmobile has a tendency to move or travel with its longitudinal axis oriented on an incline with respect to the ground level, with the front end higher than the back end. In this position, there is less weight being placed on the skis, and thus the snowmobile is more difficult to maneuver and turn. In addition, this can provide vision problems as the rider's line of sight to the ground surface is obstructed causing him/her to have to look a greater distance past the front of the snowmobile to see the ground surface. This problem is made worse at night time since it causes the headlight of the snowmobile to point uselessly toward the sky instead of on the ground in front of the moving machine."}
{"text": "Optical brighteners are a class of fluorescent materials that absorb light in the ultraviolet regions of the electromagnetic spectrum (e.g., less than 400 nm) and re-emit light in the violet and blue regions of the electromagnetic spectrum (e.g., greater than 400 nm). Optical brighteners are used in a variety of applications to enhance the color or appearance of materials, such as paper and fabrics. Generally, optical brighteners are used to increase the perceived whiteness of such materials by increasing the overall amount of blue light emanating from the material (i.e., reflected from the material and emitted by substances deposited on the material).\nWhile optical brighteners have been used to treat a variety of materials, some commercially-available optical brighteners have proven difficult to effectively deposit on certain substrates. For example, some commercially-available optical brighteners do not effectively deposit onto polyester-based fabrics. When such optical brighteners are used in laundry care compositions, polyester-based fabrics treated with the compositions typically show no deposition of the optical brightener and, consequently, no improvement in their whiteness index.\nAccordingly, a need still remains for optical brighteners that exhibit effective whitening on a variety of substrates and for compositions (e.g., laundry care compositions) comprising such optical brighteners. This application seeks to provide such optical brighteners."}
{"text": "This application is related to Japanese application No. 2000-222703 filed on Jul. 24, 2000, whose priority is claimed under 35 USC xc2xa7 119, the disclosure of which is incorporated by reference in its entirety.\n1. Field of the Invention\nThe present invention relates to a method of manufacturing a semiconductor device. More particularly, the invention relates to processing of a barrier film at a contact hole etching step (etching step for forming a contact hole).\n2. Description of the Related Art\nIn recent years, a semiconductor device has been made minute and multifunctional. With this, a phenomenon in which functions of various functional components formed in the semiconductor device are deteriorated by harmful elements, molecules, and the like, has become remarkable. As measures against this, there has been adopted a method for protecting the functional components against the harmful elements, molecules and the like by various protection films (barrier films). As constituent materials of the barrier films, Al2O3 (alumina), TiO2, TiN, TiON, Ta2O5, TaSiN, SiN and the like can be enumerated. One of these barrier films blocks the penetration of a harmful material by the dense film, and another takes a harmful material in the film to trap it. However, there is a problem that the barrier films prevent etching in a subsequent processing step or produce reactive secondary product materials (deposited materials) during the etching.\nHereinafter, an example of the prior art will be described with reference to FIGS. 3A to 3C.\nFirst, a resist pattern 26 for contact hole formation is formed in a state where an NSG film 22, a BPSG film 23, an alumina film (barrier film) 24, and an NSG film 25 are deposited over a silicone (Si) substrate 21 on which various functional components (not shown) are formed (FIG. 3A). In this case, what is barriered by the alumina film 24 is not limited to a semiconductor component formed just on the Si substrate 21, but all functional components under the alumina film 24 are barriered.\nConventionally, in a subsequent contact hole etching, the alumina film and a Si oxide films (NSG film and BPSG film) are dry etched at the same time under normal conditions for etching the Si oxide films, for example, by a gas including a fluorine based gas in plasma etching such as RIE. As an example of existing techniques, although the alumina film is not used for a barrier, an example of Japanese Unexamined Patent Publication No. Hei. 8 (1996)-31932 discloses a similar step. In this case, etching is stopped in the alumina film, or a reactive secondary product material (deposited material) 27 produced from alumina is deposited on the bottom of a contact hole (FIG. 3B).\nThe deposited material prevents electrical contact between the Si substrate and a metal wiring of an upper layer, and must be removed before the wiring is formed in the contact hole. As a conventional method, the deposited material is removed by a treatment using chemicals (for example, EKC270 of EKC Technology Inc., etc.) capable of dissolving the deposited material, or by physically impacting inert molecules of Ar or the like, ions, radicals or the like to the deposited material to make scattering (called sputter etching, inverse sputtering, or the like).\nHowever, the conventional methods have problems. That is, both the deposited material removal by the chemicals and the deposited material removal by the sputter etching have an etching rate with respect to the Si substrate of the bottom of the contact hole (see Table 1). Thus, there arise problems that a damage 28 is caused (FIG. 3C) at the bottom of the contact hole (surface of the Si substrate), and the contact resistance becomes unstable from fluctuation in the removable state of the deposited material.\nAs another example of the prior art, as in an embodiment of Japanese Unexamined Patent Publication No. Hei. 7 (1995)-273112, a method is also conceivable in which etching of the alumina film is carried out by changing a gas and by a gas containing chlorine. In the etching by the gas containing chlorine, the etching rate of the alumina film is fast, and an etching stop in the alumina film does not occur. However, the mechanism in which the deposited material is produced is not changed, and since a similar treatment for the deposited material becomes necessary, the superiority is not high.\nIncidentally, in the foregoing description of the prior art and the problems, the same is the case with a barrier film (film made of TiO2, TiN, TiON, Ta2O5, TaSiN, or SiN) other than the alumina film.\nAs described above, in the etching of the contact hole containing the barrier film, it is necessary to obtain contact hole etching conditions under which the deposited material is not produced from the barrier film, or to conceive a method in which the deposited material is effectively and certainly removed without damage.\nAccording to the present invention, there is provided a method of manufacturing a semiconductor device provided with a first insulating film and a barrier film on a conductive region and an opening portion in the first insulating film and the barrier film, the method comprising the steps of:\nforming a first opening portion in the barrier film reaching the first insulating film;\nforming a second insulating film at least on the first insulating film in the first opening portion; and\nforming a second opening portion smaller than the first opening portion and reaching the conductive region by simultaneously boring a hole into the first insulating film and the second insulating film in the first opening portion.\nOne of the features of the present invention is that the barrier film is removed more extensively than the second opening portion in advance to suppress the production of a deposited material in the opening portion at the etching for forming the second opening portion (contact hole), so that the barrier film which causes deposit in the contact hole is not exposed at the etching of the second opening portion.\nFurther, according to the present invention, the step of forming the first opening portion in the barrier film reaching the first insulating film may comprise covering the barrier film with a third insulating film, forming a resist pattern by using a mask for forming the second opening portion, forming a pattern for forming the first opening portion on the third insulating film by using the resist pattern, and forming the first opening portion in the barrier film by using the pattern for forming the first opening portion.\nAccording to the above method, since the formation of the first opening portion can be made by using the second opening portion formation mask, the number of masks to be used can be decreased.\nThese and other objects of the present application will become more readily apparent from the detailed description given hereinafter. However, it should be understood that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description."}
{"text": "In many environments, it is useful to have a surface which is level regardless of the orientation of the structure supporting the surface. For example, helicopters or other vehicles including fixed-wing Unmanned Aerial Vehicles (UAVs) and rotary-wing Vertical Take-off Unmanned Aerial Vehicles (VTUAVs) often need a level surface from which they can take off and land in a wide range of sea conditions. Providing such a surface is not a problem on land-based facilities or other structures that are secured to a solid surface such as oil drilling platforms that are anchored to a sea bed, etc. However, providing such a surface can be difficult when the object on which the aircraft is to land can move. For example, helicopters attempting to take off and land on ships are currently limited to a relatively narrow range of sea conditions in order to avoid injuring the aircraft and/or crew.\nGiven these problems, there is a need for a mechanism that can be operated to maintain a surface in a desired orientation regardless of the fact that the support for such a surface is not level and/or is moving."}
{"text": "A vertical field-effect transistor (FET) has a channel perpendicular to the substrate surface, as opposed to being situated along the plane of the surface of the substrate. By using this design, it is possible to increase packing density. That is, by having the channel perpendicular to the substrate surface, vertical FETs improve the scaling limit beyond planar finFETs.\nHowever, vertical FETs are still severely challenged past the 7 nm node due to high aspect ratios, Vmax limits, and material thickness not scaling well. For example, insulator material and shared contacts formed between gate material of adjacent vertical FETs make it very difficult to scale the devices beyond the 7 nm node, basically due to material thicknesses, leakage concerns, breakdown voltage, decreased resistances and capacitance, etc. Accordingly, these constraints make it very difficult to decrease gate pitch in current vertical FET designs."}
{"text": "The present application claims priority under 35 U.S.C xc2xa7119 to Japanese Patent Application No.2000-073932 filed Mar. 13, 2000 entitled xe2x80x9cROLLING GUIDE DEVICE xe2x80x9d, No.2000-367605 filed Dec. 1, 2000 entitled xe2x80x9cROLLING GUIDE DEVICE AND DRIVE SYSTEM USING ROLLING GUIDE DEVICExe2x80x9d, and No.2001-037486 filed Feb. 14, 2001 entitled xe2x80x9cROLLING GUIDE DEVICE AND DRIVE SYSTEM USING ROLLING GUIDE DEVICExe2x80x9d. The contents of that application are incorporated herein by reference in their entirety.\n1. Field of the Invention\nThe present invention relates to a rolling guide device in which a movable rail is made slidable with respect to a track rail and also relates to a drive system using such a rolling guide device.\n2. Description of the Related Art\nAs a rolling guide device in which a movable rail is made slidable with respect to a track rail, there has been known a slide rail such as shown in FIG. 25 (see Japanese Utility Model Publication No. SHO 62-8765). Such slide rail comprises a track rail member 1 having an opened recess 1a (i.e. having substantially -shaped (box-shaped) cross section) formed by both inner side surfaces 1b, 1b and a bottom surface 1c and a movable rail member 2 which is supported between both the inner side surfaces 1b, 1b of the track rail member 1 to be movable in the longitudinal direction thereof. The movable rail member 2 also has an opened recess 2a (i.e. having substantially -shaped (box-shaped) cross section).\nThe track rail member 1 and the movable rail member 2 have substantially the same longitudinal length. The inner side surfaces 1b of the track rail member 1 are formed with ball rolling grooves, respectively, along which balls roll in the longitudinal direction thereof, and outer side surfaces 2b of the movable rail member 2 are also formed with loaded ball rolling grooves, respectively, so as to extend in the longitudinal direction thereof and to oppose the ball rolling grooves formed to the track rail member 1.\nA number of balls 3 are arranged and housed between these ball rolling grooves and loaded ball rolling grooves, and these balls 3 are held by a cage 4 to be rotatable and slidable. When the movable rail member 2 slides with respect to the track rail member 1 in the longitudinal direction thereof, these balls 3 roll and, hence, the slide rail becomes smoothly expandable or contractive.\nFurther, though not shown, there is also known a cam-follower type drawer device of a structure that is movable and track rails are both provided with wheels so that the movable rail is drawn with respect to the track rail, as a rolling guide device in which a movable rail is slidable with respect to a track rail.\nHowever, in the conventional slide rail such as mentioned above, a number of balls 3 disposed and arranged between the track rail member 1 and the movable rail member 2 do not completely perform the rolling motion and will roll with a slight sliding motion. In the conventional slide rail, because the balls 3 do not circulate and only reciprocally move along the loaded rolling passage between the ball rolling grooves and the loaded ball rolling grooves, if the balls 3 slide, the cage 4 supporting (bearing) the balls 3 would be displaced from the initial position. As a result, in spite of the fact that an effective stroke of the movable rail member 2 is not achieved, the cage 4 collides with a stopper 5 of the track rail member 1 and, hence, such effective stroke could not be obtained. In this case, when it is required to slide the movable rail member 2 with the cage 4 colliding with the stopper 5, the movable rail member 2 will slide with the balls 3 being slipped, and accordingly, a large force is required to move the movable rail member 2.\nFurthermore, in the conventional structure of the slide rail, in order to obtain a large stroke of the movable rail member 2, it is necessary for the movable rail member 2 to once come off from a portion at which the balls 3 exist and then to be engaged with that portion at which the balls 3 exist. That is, in the case where the movable rail member 2 comes off from the portion at which the balls 3 exist, for example, the movable rail member 2 which has been loaded with ten (10) balls 3 is loaded with, for example, six (6) balls 3, and hence the ability for bearing moment load, radial load and thrust load is deteriorated.\nMoreover, with the cam-follower type drawer device, because the wheels generally have backlash or looseness, the movable rail member 2 does not smoothly slide, and furthermore, because the wheel has a cylindrical structure, a direction along which a load is received is determined, and hence, the thrust load cannot be received.\nAn object of the present invention is to substantially eliminate defects or drawbacks encountered in the prior art mentioned above and to provide a rolling guide device capable of taking a large expansion/contraction stroke and sufficiently bearing moment load, radial load and thrust load at any expanded (contracted) attitude and also provide a drive system incorporated with such a rolling guide device.\nThis and other objects can be achieved according to the present invention by providing, in one aspect, a rolling guide device comprising:\na track rail formed with a rolling member rolling surface extending along a longitudinal direction thereof;\na movable rail formed with a loaded rolling member rolling surface extending along a longitudinal direction thereof so as to oppose the rolling member rolling surface of the track rail;\na track rail side rolling member circulation passage formed to the track rail so as to circulate the rolling members rolling between the track rail and the movable rail;\na movable rail side rolling member circulation passage formed to the movable rail so as to circulate the rolling members rolling between the track rail and the movable rail; and a number of rolling members disposed and arranged in the track rail side rolling member circulation passage and the movable rail side rolling member circulation passage.\nAccording to this aspect, when the movable rail slides with respect to the track rail, the rolling members arranged between the track rail and the movable rail endlessly circulate in the track rail side rolling member circulation passage and the movable rail side rolling member circulation passage while rolling therealong. As mentioned, because the rolling members circulate in the endless manner, even if the rolling member slides during the rolling motion, there is no causing of a case that a cage is shifted from the initial position as in the conventional structure, and hence, a large expansion (contraction) stroke is obtainable. Furthermore, in an optional expanded (contracted) attitude, there remains a considerable distance between the track rail side rolling member circulation passage and the movable rail side rolling member circulation passage, so that a rolling guide device, which can bear even the moment load, can be realized.\nFurther, when the movable rail slides and its stroke is made large, the considerable distance corresponding to this stroke is made short and the capability of bearing the moment load is reduced. However, according to the present invention, the movable rail does not come off from the balls, so that the capability of bearing the moment load is not extremely reduced. Moreover, because the movable rail does not come off from the balls and the number of the rolling members supported at an optional expansion (contraction) attitude is not changed, in contrast to the conventional slide rail, there can be provided a rolling guide device bearing the constant radial load and thrust load.\nIn the above aspect, the following preferred embodiments or examples may be provided with advantageous functions and effects thereof.\nThe track rail side rolling member circulation passage is formed in one longitudinal end side of the track rail and the movable rail side rolling member circulation passage is formed in one longitudinal end side, opposing that one end side of the track rail, of the movable rail.\nAccordingly, the distance between the track rail side rolling member circulation passage and the movable rail side rolling member circulation passage can be made large, so that a rolling guide device bearing the large moment load can be provided.\nFurthermore, the track rail has an opened recess having a-shaped section and has inside surfaces to which the rolling member rolling surfaces are formed, the movable rail is fitted into the recess of the track rail, and the movable rail has outside surfaces to which the loaded rolling member rolling surfaces are formed so as to oppose the rolling member rolling surfaces formed to the track rail inside surfaces.\nThus, various kinds of loads including radial load, thrust load and moment load can be supported in a balanced condition.\nThe track rail side rolling member circulation passage is provided with a rolling member return passage substantially parallel to the rolling member rolling surface and a rolling direction changing passage communicating the rolling member rolling surface and the rolling member return passage, the movable rail side rolling member circulation passage is provided with a rolling member return passage substantially parallel to the loaded rolling member rolling surface and a rolling direction changing passage communicating the rolling member rolling surface and the rolling member return passage, the rolling direction changing passages of the track rail side rolling member circulation passage and the movable rail side rolling member circulation passage are formed to a deflector which is formed independently from a track rail body and a movable rail body, and the deflector is fitted to holes formed in the track rail body and movable rail body from the side portions thereof.\nAccording to this embodiment, the rolling direction changing passages can be easily formed in the long track rail and movable rail.\nThe return passages are drilled in the track rail body and the movable rail body from the longitudinal end portions thereof.\nAccording to this embodiment, the return passages can be easily formed in the long track rail and movable rail.\nThe deflector is composed of a plurality of sections splitting along the rolling direction changing passages.\nAccordingly, the rolling direction changing passages having a complicated structure may be easily formed in the deflector.\nThe deflector is made of a synthetic resin.\nAccordingly, the rolling direction changing passages having a complicated structure may be easily formed in the deflector, and moreover, noise which may be generated when the rolling members roll in the rolling direction changing passages will be suppressed.\nThe above mentioned object of the present invention can be also achieved by providing, in another aspect, a drive system comprising:\na track rail formed with a rolling member rolling surface extending along a longitudinal direction thereof;\na movable rail formed with a loaded rolling member rolling surface extending along a longitudinal direction thereof so as to oppose the rolling member rolling surface of the track rail;\na track rail side rolling member circulation passage formed on the track rail so as to circulate the rolling members rolling between the track rail and the movable rail;\na movable rail side rolling member circulation passage formed on the movable rail so as to circulate the rolling members rolling between the track rail and the movable rail;\na number of rolling members disposed and arranged in the track rail side rolling member circulation passage and the movable rail side rolling member circulation passage; and\na linear motor means having a primary side mounted to either one of the track rail and the movable rail and a secondary side mounted to another one of the track rail and the movable rail.\nAccording this aspect, the expansion (contraction) stroke can be made large and the moment load, the radial load and the thrust load can be sufficiently supported at an optional attitude of the system. Furthermore, because the linear motors are incorporated between the track rail and the movable rail, the use of the ball screw or like can be eliminated, thus moving the movable rail at high speed with less noise. Moreover, because it is not necessary to provide a space for a rotary motor, the drive system can be made thin and compact.\nAccording to preferred embodiments or examples of this aspect, the following advantageous functions and effects may be attained.\nThe track rail side rolling member circulation passage is formed in one longitudinal end side of the track rail and the movable rail side rolling member circulation passage is formed in one longitudinal end side, opposing that one end side of the track rail, of the movable rail, and the linear motor means comprises first and second linear motors, the first linear motor having a primary side mounted to a portion near the track rail side rolling member circulation passage of the track rail, the second linear motor having a secondary side mounted to the track rail along the longitudinal direction thereof so as to be continuous to the primary side of the first linear motor, and the second linear motor having a primary side mounted to a portion near the movable rail side rolling member circulation passage of the movable rail, the first linear motor having a secondary side mounted to the movable rail along the longitudinal direction thereof so as to be continuous to the primary side of the second linear motor.\nAccording to this embodiment, because two sets of linear motors are incorporated in the drive system, the thrust force can be made two times (twice), and the excitation is averaged to thereby make smooth the movement of the movable rail. Furthermore, the first linear motor has a primary side mounted to a portion near the track rail side rolling member circulation passage of the track rail and the second linear motor has a primary side mounted to a portion near the movable rail side rolling member circulation passage of the movable rail, so that the thrust force can be generated at substantially the same positions of the movable side rolling member circulation passage and the track rail side rolling member circulation passage, regardless of the stroke of the movable rail. Therefore, even if a pitching or yawing moment is applied to the movable rail, the thrust force can be stably applied to the movable rail.\nThe first and second linear motors may be composed of linear induction motors or linear pulse motors such that the secondary sides thereof are opposed to each other.\nFor example, in a case where linear D.C. motors are used, two sets of linear motors are disposed in back-to-back arrangement and a distance between the secondary side magnets is short, an alternating magnetic field may be generated between the magnets. However, according to this embodiment of the present invention, because the linear induction motors or linear pulse motors are used without using the magnets, there is no fear of causing any alternating magnetic filed. However, a linear D.C. motor may be utilized as far as a relatively large distance between the secondary sides of the linear D.C. motors can be taken so as not to influence from each other.\nThe nature and further characteristic features of the present invention may be made clear from the following descriptions made with reference to the accompanying drawings."}
{"text": "Many modern vehicles (e.g., automobiles or aerial vehicles) or mobile computer devices, such as smartphones, tablet computers or other general or special purpose machines, include one or more components for determining information regarding positions, orientations, velocities or accelerations of such vehicles or devices, and for providing navigational advice or instructions based on such information. For example, some vehicles and computer devices include Global Positioning System (or “GPS”) transceivers for determining positions using data received from one or more orbiting satellites, or cellular telephone equipment configured to estimate (e.g., triangulate) a position using signals received from one or more cellular telephone network towers or other network sources.\nA GPS-enabled vehicle, device or other system or component may determine its position by interpreting signals that are received from multiple GPS satellites. A distance between a GPS receiver and a GPS satellite may be determined by calculating a “time of flight” between the GPS satellite and the GPS receiver for each such signal, which is assumed to travel at approximately the speed of light. Where three or more such signals are interpreted, the GPS receiver may be determined to be located at a specific point on the planet to within a certain degree of accuracy or tolerance, commonly on the order of two to ten meters.\nOccasionally, position information determined using GPS satellites and GPS-enabled equipment may be inaccurate, irrelevant or unavailable. For example, like any computer device, most GPS-enabled equipment requires an initialization period during which GPS position information obtained or determined by such equipment is unreliable. Furthermore, where an environment includes many natural or artificial obstructions, such as tree limbs, office towers, mountains, walls or ceilings, the receipt of GPS signals by a GPS-enabled device may be delayed or otherwise interpreted as having arrived in an untimely manner. Moreover, the use of GPS satellites and GPS-enabled equipment necessarily relies on the viability of the GPS system as a whole. Like any computer-based system or communications system, the GPS system is at risk, at least theoretically, of hacking, spoofing or jamming, and may be taken offline or rendered unavailable for one or more legitimate or illegitimate purposes from time to time."}
{"text": "The exemplary embodiment relates to color selection. It finds particular application in connection with a system and method for classification and selection of color palettes based on an achromatic taxonomy.\nA color palette is a set of different colors, which may be displayed as an ordered or unordered sequence of swatches, one for each color in the set. The colors in the palette have typically been selected by a graphic designer, or other skilled artisan working with color, to harmonize with each other, when used in various combinations. The effect of different color palettes on the human perception of the objects to which the colors are applied, such as brochures, textiles, automobiles, and the like, is not fully understood. However, either through common usage or innate preferences, the selection of a color palette can influence a customer's perception of an object. For example a brochure advertising Mediterranean vacations which utilizes a set of muted hues may be less likely to draw the interest of a customer than one which uses the bright colors typically associated with the region. Accordingly, designers often wish to select a color palette which is associated, in human perception, with a certain achromatic entity, such as an emotion, geographical location, musical genre, seasons of the year, or the like. In the case of emotions, palettes which elicit an excited or enthusiastic emotion may be suitable for some applications, such as a flyer for a circus production, while for others, such as a food product, the designer may wish to elicit emotions of pleasantness, such as contented or warm-hearted emotions. The designer uses his or her acquired experience in selecting a suitable color palette for the particular application.\nIncreasingly, color palettes are used by those who do not have extensive experience in color selection and design. To ascertain the response of a wide variety of people to a particular color palette, the palette may be submitted for review. Users look at the palettes, e.g., displayed on a screen or printed on print media, and assign labels which are natural language expressions of the particular users' perceptions. Large collections of labelled data where millions of users create, annotate, and review color combinations are now accessible.\nThe present exemplary embodiment takes advantage of such labelled data in a system and method which enables selection and creation of palettes according to achromatic entities, such as emotions."}
{"text": "The literature is full of examples where processor and system faults or \"bugs\" were discovered long after the processors or systems were shipped to customers. It is well known that the later in the product cycle that a \"bug\" is discovered, the greater the expense to fix it. Compounding this problem is the trend towards shorter and shorter product cycles. Finally, the problem is compounded again by the trend towards tightly-coupled multiple processor computer systems. This compounding is because in such a tightly-coupled multiple processor system, it is not only necessary to discover and fix the faults in a single processor, it is also now necessary to discover and fix faults resulting from the interaction among the multiple processors.\nOne problem with implementing tightly coupled multiple processor computer systems are in exhaustively testing the interactions between and among multiple processors. For example, in a tightly coupled system, two or more processors may each have an individual high-speed level one (L1) cache, and share a slightly lower speed level two (L2) cache. This L2 cache is traditionally backed by an even larger main memory. The L1 and L2 caches are typically comprised of high speed Static Random Access Memory (SRAM), and the main memory is typically comprised of slower speed Dynamic Random Access Memory (DRAM).\nIt is necessary that the cache and memory be maintained for coherency. Thus, for example, at most only a single L1 cache of a single processor is allowed to contain a cache line corresponding to a given block of main memory. When multiple processors are reading and writing the same block in memory, a conflict arises among their cache controllers. This is conflict is typically resolved in a tightly coupled multiprocessor system with an interprocessor cache protocol communicated over an interprocessor bus. For example, a first processor may be required to reserve a cache copy of the contested block of memory. This is communicated to the other processors. However, if another (second) processor already has reserved the contested block of memory, the first processor must wait until the block is unlocked, and potentially written at least back to the L2 cache.\nDebugging a cache protocol can be quite difficult. This stems from a number of interrelated factors. First, the multiple processors are each typically operating asynchronously from each other at extremely high frequencies or rates of speed. Secondly, the L1 caches, and their cache controllers are typically operating at essentially the same speed as the processors. Third, instruction cache misses for test instruction sequences can delay instruction execution by relatively long, somewhat variable, periods of time. There are a number of reasons for this later problem. One reason is it may be possible to retrieve a cache line of instructions from L1 cache or from L2 cache, or it may be necessary to load the cache line from slower main memory. The DRAM comprising the main memory typically operates quite a bit slower than the processor (and L1 cache). Another problem is that the time it takes to fetch a block of instructions from the main memory may vary slightly. There are a number of causes of this. First, accessing different addresses in the DRAM may take slightly different times. This is partly because of differing signal path lengths. Secondly, different memory banks may have slightly different timing. This is true, even when the specifications for the memories are equivalent. This is particularly true, when the memories are self-timed. This problem may be accentuated when multiple processors or multiple memories share a common memory access bus, where the actions of one processor or memory may lock out, and stall, another processor or memory. Note also that asynchronous Input/Output (I/O) operations to memory can have seemingly random effects on timing.\nDespite the problems described above, in order to effectively test the interaction among multiple processors, it is preferable to exhaustively test each set of possible combinations. In the case of a cache protocol as described above, it is preferable to exhaustively test each possible set of cache states and cache state transitions. It is also preferable to be able to detect and record state changes at a lower level than that available to a user program.\nIn order to test the interactions among multiple processors, the various combinations of states and state transitions should be tested. This can be done by executing programs simultaneously on each of the processors. Varying the time when each processor executes its program can test the different combinations. Unfortunately, there is no mechanism in the prior art to accurately exhaustively vary the times when each processor executes its program. This is partly due to the processor instruction timing variations described above. The result is that timing windows often arise where particular state and state transition interactions are not tested.\nOne solution to this problem is to increase the number of tests run and the number of test cycles run. This increases the chances of uncovering faults, but does not guarantee exhaustive fault coverage.\nAnother set of prior art solutions is to try to control more closely the timing between executions of programs by the multiple processors. One such solution is to use NOP instructions to delay execution. The larger the number of NOP instructions executed, the longer the delay. However, NOP instructions are typically executed out of blocks of instructions held in cache lines. Each time execution crosses a cache line boundary, there is a potential for a cache miss, resulting in retrieving the cache line from slower memory. There is also a potential at that point that execution may be delayed for one or more cycles due to memory bus contention. Each of these potential delays introduces a potential window that did not get tested utilizing this set of solutions. Note also that virtual memory program activity must also be accounted for.\nAnother problem that arises is that it is often hard to distinguish states and state transitions from a programmer's view of a processor. This is partly because there is much that is not visible at this level. States and state transitions must therefore be assumed from visible programmer model level changes in the processor. This problem of distinguishing state and state transitions is a particular problem when the states and state transitions are cache states and state transitions during interaction testing among multiple processors.\nOne prior art solution to determining machine states and state transitions is through the use of SCAN. Using SCAN, a known pattern of states can be loaded into a processor. The processor then executes one or two instructions. The states of the various memory elements in the processor are then unloaded from the processor and compared with their expected values. This type of functional testing is becoming common for high-end microprocessors. Unfortunately, it does not lend itself to exhaustively testing the interactions among multiple processors. One reason for this is that a processor under the control of SCAN typically only executes for one or two instruction cycles, before the SCAN latches are unloaded, and another set of values loaded. The result of this is that SCAN is extremely slow, especially in comparison to the speed of modern processors. This significantly reduces the amount of testing that can be realistically done with SCAN. Secondly, there is no readily apparent mechanism available to test multiple processors at the same time, and more importantly to vary the start times of each of the multiple processors being tested together.\nIn the past, it has been sometimes been possible to run enough signals out of a processor that the states and state transitions being tested can be monitored by test equipment. One problem with this method of testing is that it is a manual and error prone process. Just as important, this method is fast becoming less and less possible as more and more functionality is embedded on single chips. Pin-count has become a major concern, and it has become increasingly unlikely that precious external pins can be dedicated for the sort of interprocessor state testing described above.\nTestability, and thus reliability through earlier fault detection would be significantly increased in tightly coupled multiprocessor systems if the interactions among multiple processors could be accurately exhaustively tested, with the guarantee that no timing windows were inadvertently left untested. This testability would be further enhanced by a mechanism for recording states and state transitions over a series of clock cycles for each of the processors being tested."}
{"text": "A Fixed Wireless Loop (FWL) is a system that supports access of wireline telecommunications devices to the public wireline network via wireless communications. A radio transceiver transmits and receives communication signals to and from a remote transceiver located at a subscriber's home or office. Thus, a FWL eliminates the need for copper wire transmission between service providers and their respective subscribers.\nA FWL is typically provided utilizing a FWL terminal at the subscriber's home or office. Mobility has been added to the FWL via a cradled handset. While the handset, or hand-held transceiver, is in the cradle, the combined unit functions as a FWL terminal providing telephone service over the subscriber's telephone wiring. When the hand-held transceiver is removed from the cradle, the hand-held transceiver acts as a mobile unit allowing the subscriber to roam throughout the house, or the FWL home sector. This type of FWL terminal serves the market of telecommunications users not needing mobile wireless services on a continuous basis, but would like to have easy access to such services when needed. Known FWL terminals as described above include the CelDock manufactured by Telular and the IntelliGo manufactured by GTE.\nSuch devices, however, render the wireline phones connected to the FWL terminal inoperative when the hand-held transceiver is removed from the cradle. Furthermore, the mobility of the hand-held transceiver is limited to the FWL sector surrounding the subscriber's home or office.\nThus, there exists a need for a FWL terminal that supports wireline telecommunications services when the hand-held transceiver is removed from the cradle. There also exists a need for providing hand-off capability between the FWL hand-held transceiver and the public wireless network."}
{"text": "1. Field of the Invention\nThis invention relates generally to a tensioner for tensioning a power transmission belt of a belt drive system. Particularly, it relates to a tensioner having a linear pulley movement. Specifically, this invention relates to a linear movement tensioner whose movement is mechanically damped and travel limited with included applicability to power transmission belt drive systems having a unitary device performing both the engine starting function and the electrical power generation function, such as a motor/generator sometimes referred to as a Gen-Star.\n2. Description of the Prior Art\nInternal combustion engines, or power plants, commonly use power transmission belt drive systems to tap power from the engine's crankshaft and deliver it to one or more various engine auxiliaries or accessories. In automotive applications, these accessories include power steering pumps, water pumps, air conditioning compressors, fuel pumps, and alternators. Historically, such engines have had the main power takeoff point at the crankshaft protruding from the rear of the engine to which is attached the drive train for driving the wheels to move the automobile. The accessories are driven from a pulley attached to the front of the crankshaft. Each accessory is equipped with a pulley. All of the pulleys are in mechanical communication via one or more power transmission belts trained about them. Some method of tensioning each power transmission belt is provided. The power transmission belt, the pulleys, and devices accomplishing belt tensioning form the accessory belt drive system.\nEarlier systems included multiple v-belts. Commonly, each belt was tensioned by manual adjustment and fixing of the position of one accessory or idler per belt. These are referred to as locked-center belt drives, because there is no provision for automatic movement of any of the pulleys to accommodate varying condition of the belt or of the drive as a whole. If the belt should stretch or otherwise lengthen, the tension upon the belt would lessen. Further, for proper operation of the belt drive system, the tension of the belt must be set high enough to accommodate the worst case condition. Such worst case conditions can be the result of extremes of temperature, engine operation, or accessory operation.\nThere has been interest in making the volume, of the engine compartments of automobiles, smaller. To accommodate the smaller compartments, various aspects of the engines have become smaller, including the accessory belt drive systems. This has been accomplished, at least in part, by reducing the number of belts employed. As each belt is removed, and the number of layers extending from the front of the engine is thereby removed, the total distance the belt drive system extends from the front of the engine is reduced. Ultimately, this has resulted in the use of a single serpentine belt for many applications. A serpentine belt is so named because of the way it snakes around the various pulleys in a series of bends, both forward and backward. A v-ribbed or Micro-V (a registered trademark of The Gates Rubber Company) belt is most suited to serpentine applications.\nThe limitations of the locked-center approach to belt tensioning are exacerbated in serpentine applications. Accordingly, most modern serpentine belt drives include an automatic tensioner whereby the changing conditions of the belt drive system can be better accommodated. In basic form, an automatic tensioner has a framework or attachment point, which attaches directly to the cylinder block of the engine, or indirectly to some point on the vehicle that is stationary with respect to the engine of the vehicle, and a pulley, which presses upon the belt in the plane of rotation of the belt drive system. A moveable member or connective portion extends between the framework and the pulley to provide pressure upon the belt, via the pulley. The pressure acts to lengthen the distance about which the belt is trained and thereby causes the belt to be in tension. Various techniques and geometries have been employed to provide the biasing force. Commonly, a resilient member, such as a steel spring acts to force the moveable member in motion which results in the pulley tending to move in a direction toward a surface of the belt which, in turn, tends to increase tension upon the belt.\nA tensioner with only these elements provides a somewhat constant force upon the surface of the belt when the system is in a resting state (i.e., the pulleys are not rotating). Dimensional instability, of the drive system caused by time, temperature, or manufacturing variation is accommodated fairly well through the action of the resilient member, at least to the limits of the resilient member and geometry of the tensioner. Thus, the tension upon the belt remains relatively constant, when the system is at rest, even though the belt may have stretched or the engine may be hot or cold. However, a tensioner with only these elements may not maintain appropriate tension upon the belt for all operating conditions of the system.\nAn operating power transmission belt drive system typically oscillates due to the influences of torsional vibration or other angular acceleration of the crankshaft or accessories, the influences of unbalanced conditions, or other influences. Torsional vibration of the crankshaft occurs, in part, as a result of the distinct impulses delivered to the crankshaft through the combustion cycles of each cylinder and piston combination. The oscillations lead to vibration of the belt. This, in turn, leads to vibration of the moveable portions of the tensioner. Momentum then builds in those moveable portions modifying the force the pulley exerts upon the belt surface and the tension upon the belt. The changing tension upon the belt can cause unacceptable performance for the belt drive system. In one instance, issues of short-term performance, such as where the belt of the belt drive system slips excessively limiting the system's efficiency or power transmission capability, or is excessively noisy due to slippage or otherwise, can arise. In another instance, the amount of tension necessarily applied to the belt, to have acceptable performance on the short-term, leads to long-term issues such as premature failure of one or more components of the system, including the belt, or one or more accessories.\nTo accommodate these issues and thus improve the performance of tensioners, damping devices have been included in tensioners. Early damped tensioners have included symmetrical damping where movement of the moveable portions of the tensioners are damped approximately equally whether the instantaneous movement is in the direction tending to increase tension upon the belt or in the direction tending to decrease tension upon the belt. Damping combines with the forces supplied by the resilient member to result in a modified biasing, at the pulley/belt interface.\nOther tensioners have utilized asymmetrical damping. Commonly, such tensioners are damped such that the damping upon the moveable portion is minimal when the tensioner is moving in the belt tensioning direction and maximal when moving in the belt loosening direction. In one approach, a shoe is biased against a race at an angle different from normal to the surface of the race. As a result, the relative movement of the shoe and race in one direction tends to lift the shoe from the race. This reduces the pressure at their interface, reduces the friction that gives rise to the damping, and thereby reduces the damping. The other direction tends to wedge the shoe against the race and increase the damping. An example is described in U.S. Pat. No. 5,964,674, to Serkh et al. These have involved the use of tensioners having a single pulley biased against a surface of the belt, to supply tension. Further, the biasing against the belt has been solely relative to the engine block.\nU.S. Pat. No. 4,416,647 to White, Jr. discloses the use of tensioners with two pulleys pressing upon the power transmission belt. The '647 patent states that the approach is useful for tensioning a system with a cyclic load such as an air conditioning compressor. One, of the pulleys, presses upon a span of the power transmission belt immediately upstream of the cyclic load. While, the other pulley presses upon the power transmission belt immediately downstream of the cyclic load. In one embodiment, the two pulleys are fixed relative to each other upon a angled member that can pivot about its apex. The assembly is pressed toward the power transmission belt to supply static tension in a locked-center fashion. The pivot is said to accommodate dynamic tension. Static tension is the result of the force applied to the power transmission belt by the tensioner in the belt tensioning direction with the effect of tending to lengthen the distance the power transmission belt is forced to travel about the pulleys of the system. If it were assumed that each of the pulleys of the system is allowed to rotate freely, tension on every span would be the same and at static tension. Dynamic tension is the tension over the length of the power transmission belt that is the result of static tension as altered by the influences of torque upon each of the pulleys and various imbalances of the system. As an additional result, each span tends to be under differing tension.\nPower transmission belt tensioners used in conjunction with accessory belt drive systems for internal combustion engines are known with pulley movements that are both arcuate and linear. Those of the above discussion have all used arcuate movements. Certain applications benefit from linear movement tensioners (i.e., those where the idler pulley contacting the power transmission belt moves in a substantially straight path).\nAn example of a linear movement tensioner can be found in U.S. Pat. No. 4,634,408 to Foster. The '408 patent discloses the construction of a relatively complex tensioner that incorporates a plurality of springs about a hydraulic damping mechanism. The tension upon the belt is adjustable by bolts that define the level of compression applied to the springs. The carrier, upon which the tensioning pulley is mounted, engages the track, that defines the carriers two degrees of freedom of movement, with a ball and track bearing structure. All of this, has inherent limits of travel, as with all mechanical structures. However, there is no intentional travel limitation disclosed. Further, there is no disclosure of travel limitation structure or technique to accommodate a controlled deceleration if the tensioner reaches the limit of it travel. It is common in the art, of tensioners for belt drive systems, for tensioners to be designed to have adequate travel in their movements to never reach the limit of their travel in normal belt system operation. The '408 patent does not suggest that the disclosed tensioner is contemplated to be any different. The inclusion of hydraulics creates the concern of leaking hydraulic fluid. The hydraulic mechanism disclosed is apparently not compact. Thus, the tensioner appears to be relatively large. The relatively large number of parts adds to cost and complexity of construction. Further, there is no provision for compensating for the parasitic torque across the carrier that supports the idler pulley, which tends to twist that carrier, other than the effectiveness of the ball and track bearing structure.\nTraditionally, an electric starter motor is provided to spin the crankshaft of the engine so that combustion may be initiated and the engine will begin to run. The starter motor is located near the rear of the engine and is adapted to intermittently engage the rear portion of the crankshaft through a gear train.\nCurrently, there is increasing pressure to reduce emissions and increase fuel economy by lowering the weight of the automobile and reducing the number of under-the-hood components. An approach taken toward these goals involves combining the function of the starter motor and the function of the alternator into a single device, a motor/generator or a Gen-Star. Also toward the goal of increasing fuel economy, the Gen-Star promotes the use of a feature called “stop-in-idle”. This feature is where the engine is allowed to die when it would ordinarily idle, then be restarted when the automobile is expected to resume motion. This feature substantially increases the demands placed upon accessory belt drives. In application, the motor/generator is placed in mechanical communication with the crankshaft via the accessory belt drive. The motor/generator and associated accessory belt drive system tends to be placed at the front of the engine. However, placing these systems at other locations, including the rear of the engine is envisioned.\nThe advent of Gen-Star systems causes the designer, of power transmission belt drive systems, to face substantial new challenges, above mere oscillatory loads. A significant challenge, among these, has been to develop a tensioning system that results in acceptable performance, by an accessory belt drive that includes this new device, which not only offers substantial load and rotational inertia, but also adds large driving torque into the accessory belt drive. Further, it provides this large driving torque on an intermittent basis.\nIt has been discovered that for various layouts of belt drive systems incorporating a linear movement tensioner and a motor/generator that the inclusion of damping, especially asymmetric damping, can positively affect both short-term and long-term performance of such a belt drive system. It has also been discovered that placing specific limits on the travel of the tensioner in the belt loosening direction can positively affect both short-term and long-term performance of such a belt drive system. This can be true even where a system does not include a motor/generator. However, deleterious side effects caused by such a travel limitation include undesirable belt and other system noise, when the limits are reached.\nU.S. Pat. No. 4,758,208 to Bartos et al., discloses the use of two arms with arcuate movements, each carrying a pulley. The arms are mounted with pivot points that correspond with the shaft of a motor/generator. The two arms are biased toward each other by a spring. The tensioner also includes having the motor/generator mounted in a limited rotatable fashion such that the housing is allowed to rotate a few degrees in reaction to whether the motor/generator mode is of operating as a starter or an alternator. This reactive movement operates a pair of latches, which alternatingly lock one or the other of the two arms against movement, depending upon mode. In this manner, the arm associated with the power transmission belt span terminating at the motor/generator pulley with the greatest tension, as a result of the motor/generator mode, is locked in place. The free arm then supplies tension to the power transmission belt system. This tensioner is apparently complex, requires special mounting of the motor/generator having moving parts subject to wear, and is not flexible in its application. Further, the '208 patent does not disclose the contemplation of adding damping to the movement of either pulley to enhance system performance. Also, it does not disclose a travel limiting structure that is either resilient, cooperates with a damping mechanism, or suggests an attempt at overcoming the previously described deleterious side effects.\nInstallation of a linear tensioner having specific travel limits can also be enhanced by making those specific limits adjustable on a case-by-case basis and by providing a release mechanism that can be selectively released upon installation of the belt. Further, tension decay characteristics of the belt can be better accommodated where the adjusters providing the adjustable specific limits can be preset to provide an initial locked-center tensioner adjustment.\nAccordingly, there is a continuing need for a linear movement tensioner that is travel limited while avoiding deleterious side effects in operation for those applications that can benefit therefrom. An accessory power transmission belt drive system having a motor/generator is one such application.\nThere is a continuing need for a linear tensioner having travel limits that are adjustable, or can provide an initial locked-center tension, or both. There is also a continuing need for a linear tensioner having a release mechanism.\nFurther, there remains the need for a tensioner and system, for use in conjunction with a motor/generator that provides, at once, adequate short-term performance, adequate long-term performance, optimizes the width of the belt that may be used for any given application, contains cost and complexity, is flexible according to the motor/generator systems to which it can be applied and offers a linear movement."}
{"text": "The present application claims priority under 35 U.S.C. xc2xa7 119 to German patent applications No. 19830430.7 filed Jul. 8, 1998, and No. 19903126.6 filed Jan. 27, 1999. Both priority applications are entirely incorporated herein by reference.\nThe present invention relates to compounds of the formula I \nwherein A1, A2, R1, R2, R3, X and n are as defined below, which are valuable pharmaceutically active compounds for the therapy and prophylaxis of diseases, for example for cardiovascular diseases such as hypertension, angina pectoris, cardiac insufficiency, thromboses or atherosclerosis. The compounds of the formula I are capable of modulating the body\"\"s production of cyclic guanosine monophosphate (xe2x80x9ccGMPxe2x80x9d) and are generally suitable for the therapy and prophylaxis of diseases which are associated with a disturbed cGMP balance. The invention furthermore relates to processes for preparing compounds of the formula I, to their use for the therapy and prophylaxis of the abovementioned diseases and for preparing pharmaceuticals for this purpose, and to pharmaceutical preparations which comprise compounds of the formula I.\ncGMP is an important intracellular messenger which triggers a multitude of different effects via the modulation of cGMP-dependent protein kinases, phosphodiesterases and ion channels. Examples are the relaxation of smooth muscles, the inhibition of thrombocyte activation and the inhibition of the proliferation of smooth-muscle cells and of leukocyte adhesion. cGMP is produced by particulate and soluble guanylate cyclases as a response to a number of extracellular and intracellular stimuli. In the case of the particulate guanylate cyclases, stimulation is essentially effected by peptidic messengers, such as the atrial natriuretic peptide or the cerebral natriuretic peptide. The soluble guanylate cyclases (xe2x80x9csGCxe2x80x9d), which are cytosolic heterodimeric heme proteins, in contrast, are essentially regulated by a family of low-molecular-weight factors which are formed enzymatically. The most important stimulant is nitrogen monoxide (xe2x80x9cNOxe2x80x9d) or a closely related species. The function of other factors such as carbon monoxide or the hydroxyl radical is still largely unclear. The binding of NO to the heme with formation of a penta-coordinate heme-nitrosyl complex is being discussed as activation mechanism of the activation by NO. The associated release of the histidine which is bound in the basal state to the iron converts the enzyme into the active conformation.\nActive soluble guanylate cyclases are composed of an xcex1 and a xcex2 subunit each. Several subunit subtypes have been described which differ from one another with respect to sequence, tissue-specific distribution and expression in different development stages. The subtypes xcex11 and xcex21 are mainly expressed in brain and lung, while xcex22 is found in particular in liver and kidney. The subtype xcex12 was shown to be present in human fetal brain. The subunits referred to as xcex13 and xcex23 were isolated from human brain and are homologous to xcex11 and xcex21. More recent works indicate an xcex12i subunit which contains an insert in the catalytic domain. All subunits show great homologies in the region of the catalytic domain. The enzymes presumably contain one heme per heterodimer, which is bound via xcex21-Cys-78 and/or xcex21-His-105 and is part of the regulatory center.\nUnder pathologic conditions, the formation of guanylate-cyclase-activating factors can be reduced, or their degradation may be promoted owing to the increased occurrence of free radicals. The resulting reduced activation of the sGC leads, via a weakening of the respective cGMP-mediated cellular response, for example to an increase of the blood pressure, to platelet activation or to increased cell proliferation and cell adhesion. As a consequence, formation of endothelial dysfunction, atherosclerosis, hypertension, stable or unstable angina pectoris, thromboses, myocardial infarction, strokes or erectile dysfunction results. Pharmacological stimulation of sGC offers a possibility to normalize cGMP production and therefore makes possible the treatment and/or prevention of such disorders.\nFor the pharmacological stimulation of the sGC, use has hitherto almost exclusively been made of compounds whose activity is based on an intermediate NO release, for example organic nitrates. The drawback of this treatment is the development of tolerance and a reduction of activity, and the higher dosage which is required because of this.\nVarious sGC stimulators which do not act via NO release were described by Vesely in a series of publications. However, the compounds, most of which are hormones, plant hormones, vitamins or natural compounds such as, for example, lizard poisons predominantly only have weak effects on the cGMP formation in cell lysates. D. L. Vesely, Eur. J. Clin. Invest., vol.15, 1985, p. 258; D. L. Vesely, Biochem. Biophys. Res. Comm., vol. 88, 1979, p.1244. A stimulation of heme-free guanylate cyclase by protoporphyrin IX was demonstrated by Ignarro et al., Adv. Pharmacol., vol. 26, 1994, p. 35. Pettibone et al., Eur. J. Pharmacol., vol. 116, 1985 p. 307, described an antihypertensive action of diphenyliodonium hexafluorophosphate and attributed this to a stimulation of sGC. According to Yu et al., Brit. J. Pharmacol, vol. 114, 1995, p.1587, isoliquiritigenin, which has a relaxing action on isolated rat aortas, also activates sGC. Ko et al., Blood vol. 84, 1994, p. 4226, Yu et al., Biochem. J. vol. 306, 1995, p. 787, and Wu et al., Brit. J. Pharmacol. vol. 116, 1995, p. 1973, demonstrated a sGC-stimulating activity of 1-benzyl-3-(5-hydroxymethyl-2-furyl)indazole and demonstrated an antiproliferative and thrombocyte-inhibiting action. Pyrazoles and fused pyrazoles which exhibit a sGC-stimulating activity are described in European Patent Application No. 908,456 and German Patent Application No. 19,744,027.\nA series of 2-sulfonylaminobenzoic acid N-arylamides, the N-aryl group of which carries a thio substituent, have been mentioned in the literature. These compounds in which the N-aryl group generally carries as further substituents groups which are readily oxidizable such as, for example, two hydroxy groups being in para position with respect to one another and which in this case can be regarded as hydroquinone derivatives, are auxiliaries for the preparation of photographic materials (see, for example, Chemical Abstracts 119, 105757; 120, 41858; 123, 70224; or 126, 257007). If isolated structural elements are considered then the N-aryl group in these known compounds corresponds to the group R1xe2x80x94S(O)nxe2x80x94A1 in formula I in case A1 denotes a 1,4-phenylene residue which in positions 2 and 5 carries hydroxy groups (or oxy substituents), and the number n is 0. British patent publication No. 876,526 (Chemical Abstracts 56, 15432e) discloses 3,5-dichloro-2-methylsulfonylaminobenzoic acid N-(5-chloro-2-(4-chlorophenylmercapto)-phenyl)-amide which can be used for the protection of wool against moths. Compounds covered by British patent publication No. 876,526 correspond to compounds of the formula I if simultaneously the ring A1 which comprises the carbon atoms which carry the groups C(xe2x95x90X)xe2x80x94NHxe2x80x94 and NHxe2x80x94SO2R2, together with the residues R3, is a benzene ring which carries one to four halogen atoms from the series chlorine and bromine, R2 is (C1-C4)-alkyl, X is oxygen and the group R1xe2x80x94S(O)nxe2x80x94A1xe2x80x94 is a phenylmercaptophenyl-residue (=phenylthiophenyl-) which is substituted by halogen and/or trifluoromethyl and which can also be substituted by methyl or (C1-C4)-alkoxy, and the total number of halogen atoms and trifluoromethyl groups is greater than two. Pharmacological activities of these known 2-sulfonylaminobenzoic acid N-arylamides are not disclosed.\nSurprisingly, it has now been found that the compounds of the present invention effect a strong activation of guanylate cyclase and are therefore suitable for the therapy and prophylaxis of disorders which are associated with a low cGMP level.\nThus, the present invention relates to compounds of the formula I: \nwherein\nA1 is a divalent residue from the series phenylene, naphthylene and heteroarylene which can all be substituted by one or more identical or different substituents from the series halogen, (C1-C5)-alkyl, phenyl, tolyl, CF3, NO2, OH, xe2x80x94Oxe2x80x94(C1-C5)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl, (C1-C2)-alkylenedioxy, NH2, xe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94N((C1-C3)-alkyl)2, xe2x80x94NHxe2x80x94CHO, xe2x80x94NHxe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N((C1-C3)-alkyl)2, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, heterocyclyl, CHO, xe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94S(O)nxe2x80x94(C1-C4)-alkyl, xe2x80x94S(O)n-phenyl and xe2x80x94S(O)n-tolyl;\nthe ring A2 which comprises the carbon atoms which carry the groups C(xe2x95x90X)xe2x80x94NHxe2x80x94 and NHxe2x80x94SO2R2 is a benzene ring, a naphthalene ring, a saturated or partially unsaturated 3-membered to 7-membered carbocycle, a saturated or partially unsaturated or aromatic monocyclic 5-membered to 7-membered heterocycle which contains one or more ring heteroatoms from the series N, O and S, or a saturated or partially unsaturated or aromatic bicyclic 8-membered to 10-membered heterocycle which contains one or more ring heteroatoms from the series N, O and S;\nR1 is aryl, heterocyclyl or (C1-C18)-alkyl which can be substituted by one or more identical or different residues R4 or, if the number n in the group R1xe2x80x94S(O)nxe2x80x94 is 2, R1 can also be NR5R6 or, if the number n in the group R1xe2x80x94S(O)nxe2x80x94 is O, R1 can also be xe2x80x94CN;\nR2 is aryl, heterocyclyl, NR5R6 or (C1-C10)-alkyl which can be substituted by one or more identical or different residues R4;\nR3 denotes one or more identical or different residues from the series hydrogen, halogen, CF3, OH, xe2x80x94Oxe2x80x94(C1-C7)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C7)-alkyl, xe2x80x94O-aryl, (C1-C2)-alkylenedioxy, NO2, xe2x80x94ON, NR7R8, xe2x80x94COxe2x80x94NR7R8, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, heterocyclyl, xe2x80x94S(O)nxe2x80x94(C1-C5)-alkyl and (C1-C5)-alkyl which can be substituted by one or more identical or different residues R4;\nR4 is fluorine , OH, xe2x80x94Oxe2x80x94(C1-C10)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C7)-alkyl, xe2x80x94O-aryl, xe2x80x94CN, NR7R8, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N((C1-C3)-alkyl)2, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, heterocyclyl or oxo;\nR5 is hydrogen, (C1-C10)-alkyl which can be substituted by one or more identical or different substituents R4 and/or by aryl, or is aryl, heterocyclyl, xe2x80x94COxe2x80x94NR7R8, xe2x80x94CO-aryl or xe2x80x94COxe2x80x94(C1-C10)-alkyl wherein the alkyl residue can be substituted by one or more identical or different residues R4;\nR6 independently of R5 has one of the meanings indicated for R5, or\nR5 and R6 together with the nitrogen atom to which they are bonded form a 5-membered to 8-membered saturated or partially unsaturated ring which in addition to the nitrogen atom which carries the groups R5 and R6 can contain one or more further ring heteroatoms from the series N, O and S and which can be substituted by one or more identical or different substituents from the series fluorine, (C1-C5)-alkyl, hydroxy-(C1-C3)-alkyl-, xe2x80x94(C1-C3)-alkyl-Oxe2x80x94(C1-C4)-alkyl, aryl, CF3, OH, xe2x80x94Oxe2x80x94(C1-C7)-alkyl, xe2x80x94O-aryl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C7)-alkyl, (C2-C3)-alkylenedioxy, NR7R8, CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N(C1-C3)-alkyl)2, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, CHO, xe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94S(O)nxe2x80x94(C1-C4)-alkyl, xe2x80x94S(O)nxe2x80x94NH2, xe2x80x94S(O)nxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94S(O)nxe2x80x94N((C1-C3)-alkyl)2, oxo, xe2x80x94(CH2)mxe2x80x94NH2, xe2x80x94(CH2)mxe2x80x94NHxe2x80x94(C1-C4)-alkyl and xe2x80x94(CH2)mxe2x80x94N((C1-C4)-alkyl)2 where in the substituent -(CH2)mxe2x80x94N((C1-C4)-alkyl)2 the two alkyl groups can be connected by a single bond and then together with the nitrogen atom carrying them form a 5-membered to 7-membered ring which besides that nitrogen atom and the carbon atoms can additionally contain an oxygen atom, a sulfur atom or a group NR5 as ring member;\nR7 is hydrogen or (C1-C7)-alkyl which can be substituted by one or more identical or different substituents from the series OH, xe2x80x94Oxe2x80x94(C1-C5)-alkyl, NH2, xe2x80x94NHxe2x80x94(C1-C4) alkyl and xe2x80x94N((C1-C4)-alkyl)2 where in the substituent N((C1-C4)-alkyl)2 the two alkyl groups can be connected by a single bond and then together with the nitrogen atom carrying them form a 5-membered to 7-membered ring which besides that nitrogen atom and the carbon atoms can additionally contain an oxygen atom, a sulfur atom or a group NR5 as ring member;\nR8 independently of R7 has one of the meanings of R7 or is xe2x80x94COxe2x80x94(C1-C4)-alkyl;\nxe2x80x9carylxe2x80x9d is phenyl, naphthyl or heteroaryl which can all be substituted by one or more identical or different substituent from the series halogen, (C1-C5)-alkyl, phenyl, tolyl, CF3, xe2x80x94Oxe2x80x94CF3, NO2, OH, xe2x80x94Oxe2x80x94(C1-C5)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl, (C1-C2)-alkylenedioxy, NH2, xe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94N((C1-C3)-alkyl)2, xe2x80x94NHxe2x80x94CHO, xe2x80x94NHxe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N((C1-C3)-alkyl)2, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, heterocyclyl, CHO, xe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94S(O)nxe2x80x94(C1-C4)-alkyl, xe2x80x94S(O)n-phenyl and xe2x80x94S(O)n-tolyl;\nxe2x80x9cheteroarylxe2x80x9d and xe2x80x9cheteroarylenexe2x80x9d are a residue of a monocyclic 5-membered or 6-membered aromatic heterocycle or of a bicyclic 8-membered to 10-membered aromatic heterocycle each of which contains one or more ring heteroatoms from the series N, O and S;\nxe2x80x9cheterocyclylxe2x80x9d is a residue of a monocyclic or polycyclic 5-membered to 11-membered saturated or partially unsaturated heterocycle which contains one or more ring heteroatoms from the series N, O and S and which can be substituted by one or more identical or different substituents from the series fluorine, (C1-C5)-alkyl, OH, xe2x80x94Oxe2x80x94(C1-C5)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl, NH2, xe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94N((C1-C3)-alkyl)2, xe2x80x94CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N((C1-C3)-alkyl)2, xe2x80x94COxe2x80x94OH and xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl;\nn is 0, 1 or 2;\nm is 2, 3 or4;\nX is O or NH or X is a nitrogen atom which via a single bond is attached to a ring carbon atom in the group A1 which ring carbon atom is directly adjacent to the carbon atom in A1 carrying the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94 so that the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94 together with the carbon atoms in A1 carrying it forms an anellated imidazole ring;\nin all their stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts;\nwhere, however, the compound of formula I is excluded wherein simultaneously the ring A2 which comprises the carbon atoms which carry the groups C(xe2x95x90X)xe2x80x94NHxe2x80x94 and NHxe2x80x94SO2R2 is a benzene ring which is substituted in positions 3 and 5 by chlorine, R2 is methyl, X is oxygen and R1xe2x80x94S(O)nxe2x80x94A1xe2x80x94 is a 5-chloro-2-(4-chlorophenylmercapto)-phenyl residue.\nIf groups or substituents can occur several times in the compounds of formula I such as, for example R3, R4, R5, aryl, heterocyclyl, alkyl, or the numbers n and m, they can all independently of one another have the meanings indicated and can in each case be identical or different.\nAlkyl residues can be straight-chain or branched. This also applies when they are part of other groups, for example in alkoxy groups, alkoxycarbonyl groups or amino groups, or when they are substituted. Examples of alkyl groups are methyl, ethyl, propyl, butyl, pentyl, hexyl, heptyl, octyl, nonyl, decyl, undecyl, dodecyl, tridecyl, tetradecyl, pentadecyl, hexadecyl, heptadecyl, octadecyl, the n-isomers of these residues, isopropyl, isobutyl, isopentyl, sec-butyl, tert-butyl, neopentyl, 3,3-dimethylbutyl. The term alkyl here also expressly includes unsaturated alkyl residues, I. e. alkyl residues which contain one or more double bonds and/or one or more triple bonds such as, for example, alkenyl residues and alkinyl residues. Of course, an unsaturated alkyl group contains at least two carbon atoms. Specific alkyl groups whose number of carbon atoms can vary from 1 to a given upper limit, thus also comprise unsaturated alkyl groups whose number of carbon atoms can vary from 2 to the given upper limit. Examples of such residues are the vinyl residue, the 1-propenyl residue, the 2-propenyl residue (allyl residue), the 2-butenyl residue, the 2-methyl-2-propenyl residue, the 3-methyl-2-butenyl residue, the ethinyl residue, the 2-propinyl residue (propargyl residue), the 2-butinyl residue or the 3-butinyl residue. Further, the term alkyl here also expressly includes alkyl residues in which by an internal ring closure within the alkyl group a cyclic system is formed, i. e. the term alkyl also includes saturated and partially unsaturated cycloalkyl residues and cycloalkyl-alkyl- residues (alkyl substituted by cycloalkyl). Of course, a monocyclic cycloalkyl group contains at least three carbon atoms.\nSpecific alkyl groups whose number of carbon atoms can vary from 1 to a given upper limit, thus also comprise monocyclic cycloalkyl groups whose number of carbon atoms can vary from 3 to the given upper limit, and appropriate cycloalkyl-alkyl- groups. Examples of such cycloalkyl residues are cyclopropyl, cyclobutyl, cyclopentyl, cyclohexyl, cycloheptyl and cyclooctyl which can all also be substituted by one or more identical or different (C1-C4)-alkyl residues, in particular by methyl. Examples of such substituted cycloalkyl residues are 4-methylcyclohexyl, 4-tert-butylcyclohexyl or 2,3-dimethylcyclopentyl. Furthermore, unless stated otherwise the term alkyl here also expressly includes unsubstituted alkyl residues as well as alkyl residues which are substituted by one or more, for example one, two, three or four, identical or different aryl residues. The term alkyl thus here also expressly includes arylalkyl- residues such as, for example, aryl-(C1-C4)-alkyl-, for example benzyl residues, phenylethyl residues or indanyl residues. In substituted alkyl residues, for example arylalkyl-, hydroxyalkyl- such as xe2x80x94(C1-C3)-alkyl-OH or alkoxyalkyl- such as xe2x80x94(C1-C3)-alkyl-Oxe2x80x94(C1-C4)-alkyl, the substituents can be present in any desired position.\nA saturated or partially unsaturated 3-membered to 7-membered carbocycle representing the ring A2 can be derived from the monocyclic parent systems cyclopropane, cyclobutane, cyclopentane, cyclohexane or cycloheptane. If the carbocycle is unsaturated it can contain, for example, one double bond or, in the case of a 5-membered ring, 6-membered ring or 7-membered ring, also two double bonds which can be isolated or conjugated. Double bonds can be present in any positions with respect to the groups C(xe2x95x90X)xe2x80x94NHxe2x80x94 and NHxe2x80x94SO2xe2x80x94R2, i.e., for example a double bond can also be present between the two ring carbon atoms which carry these two groups.\nUnless stated otherwise, phenyl residues, naphthyl residues and heterocyclic residues, for example heteroaryl residues, can be unsubstituted or can carry one or more, for example one, two, three or four, identical or different substituents which can be in any desired positions. Unless stated otherwise, in these residues for example those substituents can be present which are indicated as substituents of an aryl group. A preferred series of substituents that can be present in the residue aryl is formed by the substituents halogen, (C1-C5)-alkyl, phenyl, tolyl, CF3, NO2, OH, xe2x80x94Oxe2x80x94(C1-C5)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl, (C1-C2)-alkylenedioxy, NH2, xe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94N((C1-C3)-alkyl)2, xe2x80x94NHxe2x80x94CHO, xe2x80x94NHxe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N((C1-C3)-alkyl)2, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, heterocyclyl, CHO, xe2x80x94COxe2x80x94(C1-C5)-alkyl, xe2x80x94S(O)nxe2x80x94(C1-C4)-alkyl, xe2x80x94S(O)n-phenyl and xe2x80x94S(O)n-tolyl. If in compounds of the formula I nitro groups are present as substituents, in total only up to two nitro groups can be present in the molecules. If phenyl residues, phenoxy residues, benzyl residues or benzyloxy residues are present as substituents in, for example, aryl residues like phenyl residues and/or in heterocyclic residues then in these substituents the benzene ring can also be unsubstituted or substituted by one or more, for example one, two, three or four, identical or different residues, for example by residues from the series (C1-C4)-alkyl, halogen, hydroxy, (C1-C4)-alkoxy, trifluoromethyl, cyano, hydroxycarbonyl, ((C1-C4)-alkoxy)carbonyl, aminocarbonyl, nitro, amino, (C1-C4)-alkylamino, di-((C1-C4)-alkyl)amino and ((C1-C4)-alkyl)carbonylamino.\nIn monosubstituted phenyl residues the substituent can be in the 2-position, the 3-position or the 4-position, in disubstituted phenyl residues the substituents can be in 2,3-position, 2,4-position, 2,5-position, 2,6-position, 3,4-position or 3,5-position. In trisubstituted phenyl residues the substituents can be in 2,3,4-position, 2,3,5-position, 2,3,6-position, 2,4,5-position, 2,4,6-position or 3,4,5-position. Tolyl (i.e., methylphenyl) can be 2-tolyl, 3-tolyl or 4-tolyl. Naphthyl can 1-naphthyl or 2-naphthyl. In monosubstituted 1-naphthyl residues the substituent can be in the 2-position, the 3-position, the 4-position, the 5-position, the 6-position, the 7-position or the 8-position, in monosubstituted 2-naphthyl residues in the 1-position, the 3-position, the 4-position, the 5-position, the 6-position, the 7-position or the 8-position.\nThe above explanations as well as the following explanations relating to monovalent residues correspondingly apply to the divalent residues phenylene, naphthylene and heteroarylene. The free bonds via which the divalent residues are attached to the adjacent groups can be present on any ring carbon atoms. In the case of a phenylene residue they can be in 1,2-position (ortho-phenylene), 1,3-position (meta-phenylene) or 1,4-position (para-phenylene). In the case of a naphthylene residue the free bonds can be in 1,2-position (=1,2-naphthylene or 1,2-naphthalinediyl) or in 1,3-position, 1,4-position, 1,5-position, 1,6-position, 1,7-position, 1,8-position, 2,3-position, 2,6-position or 2,7-position. In the case of 5-membered ring aromatics containing one heteroatom such as, for example, thiophene or furan, the two free bonds can be in 2,3-position, 2,4-position, 2,5-position or 3,4-position. A divalent residue derived from pyridine can be a 2,3-, 2,4-, 2,5-, 2,6-, 3,4- or 3,5-pyridinediyl residue. In the case of unsymmetrical divalent residues the present invention includes all positional isomers, i.e., in the case of a 2,3-pyridinediyl residue, for example, it includes the compound in which the one adjacent group is present in the 2-position and the other adjacent group is present in the 3-position as well as the compound in which the one adjacent group is present in the 3-position and the other adjacent group is present in the 2-position.\nHeteroaryl residues, heteroarylene residues, heterocyclyl residues, heterocycles representing the ring A2 and rings which are formed by two groups bonded to a nitrogen atom together with this nitrogen atom are preferably derived from heterocycles which contain one, two, three or four identical or different ring heteroatoms, more preferably from heterocycles which contain one, two or three, in particular one or two, identical or different heteroatoms. Unless stated otherwise, the heterocycles can be monocyclic or polycyclic, for example monocyclic, bicyclic or tricyclic. Preferably they are monocyclic or bicyclic. The rings preferably are 5-membered rings, 6-membered ring or 7-membered ring. Examples of monocyclic and bicyclic heterocyclic systems from which residues occurring in the compounds of the formula I can be derived, are pyrrole, furan, thiophene, imidazole, pyrazole, 1,2,3-triazole, 1,2,4-triazole, 1,3-dioxole, 1,3-oxazole (=oxazole), ,1,2-oxazole (=isoxazole), 1,3-thiazole (=thiazole), 1,2-thiazole (=isothiazole), tetrazole, pyridine, pyridazine, pyrimidine, pyrazine, pyran, thiopyran, 1,4-dioxine, 1,2-oxazine, 1,3-oxazine, 1,4-oxazine, 1,2-thiazine, 1,3-thiazine, 1,4-thiazine, 1,2,3-triazine, 1,2,4-triazine, 1,3,5-triazine, 1,2,4,5-tetrazine, azepine, 1,2-diazepine, 1,3-diazepine, 1,4-diazepine, 1,3-oxazepine, 1,3-thiazepine, indole, benzothiophene, benzofuran, benzothiazole, benzimidazole, quinoline, isoquinoline, cinnoline, quinazoline, quinoxaline, phthalazine, thienothiophenes, 1,8-naphthyridine and other naphthyridines, pteridin, or phenothiazine, each of them in saturated form (perhydro form) or in partially unsaturated form (for example in the dihydro form or the tetrahydro form) or in maximally unsaturated form, insofar as the respective forms are known and stable.\nThus, the heterocycles which are suitable also include, for example, the saturated heterocycles pyrrolidine, piperidine, piperazine, morpholine and thiomorpholine. The degree of saturation of heterocyclic groups is indicated in their individual definitions. Unsaturated heterocycles can contain, for example, one, two or three double bonds within the ring system. 5-membered rings and 6-membered rings can in particular also be aromatic.\nThe residues derived from these heterocycles can be attached via any suitable carbon atom. Nitrogen heterocycles which can carry a hydrogen atom or a substituent on a ring nitrogen atom, for example pyrrole, imidazole, pyrrolidine, morpholine, piperazine etc., can also be attached via a ring nitrogen atom, in particular if the heterocyclic residue in question is bonded to a carbon atom. For example, a thienyl residue can be present as 2-thienyl residue or 3-thienyl residue, a furyl residue as 2-furyl residue or 3-furyl residue, a pyridyl residue as 2-pyridyl residue, 3-pyridyl residue or 4-pyridyl residue, a piperidinyl residue as 1-piperidinyl residue (i.e., piperidino residue), 2-piperidinyl residue, 3-piperidinyl residue or 4-piperidinyl residue, a (thio)morpholinyl residue as 2-(thio)morpholinyl residue, 3-(thio)morpholinyl residue or 4-(thio)morpholinyl residue (i.e., thiomorpholino residue). A residue derived from 1,3-thiazole or imidazole which is attached via a carbon atom can be attached via the 2-position, the 4-position or the 5-position.\nUnless stated otherwise the heterocyclic groups can be unsubstituted or can carry one or more, for example one, two, three, or four, identical or different substituents. Substituents in heterocycles can be present in any desired positions, for example in a 2-thienyl residue or 2-furyl residue in the 3-position and/or in the 4-position and/or in the 5-position, in a 3-thienyl residue or 3-furyl residue in the 2-position and/or in the 4-position and/or in the 5-position, in a 2-pyridyl residue in the 3-position and/or in the 4-position and/or in the 5-position and/or in the 6-position, in a 3-pyridyl residue in the 2-position and/or in the 4-position and/or in the 5-position and/or in the 6-position, in a 4-pyridyl residue in the 2-position and/or in the 3-position and/or in the 5-position and/or in the 6-position. Unless stated otherwise, for example those substituents can be present as substituents in heterocyclic groups which are indicated in the definition of the group aryl, and in the case of saturated and partially unsaturated heterocycles as further substituents also the oxo group and the thioxo group can be present. Substituents on a heterocycle as well as substituents on a carbocycle can also form a ring, i.e., to a ring system further rings can be condensed (or anellated) so that, for example, also cyclopenta-condensed, cyclohexa-condensed or benzo-condensed rings can be present. Suitable substituents on a substitutable ring nitrogen atom of a heterocycle are in particular, for example, unsubstituted (C1-C5)-alkyl residues and aryl-substituted alkyl residues, aryl residues, acyl residues such as xe2x80x94COxe2x80x94(C1-C5)-alkyl, or sulfonyl residues such as xe2x80x94SO2xe2x80x94(C1-C5)-alkyl. Suitable nitrogen heterocycles can also be present as N-oxides or as quaternary salts containing a counterion which is derived from a physiologically acceptable acid. Pyridyl residues, for example, can be present as pyridine-N-oxides.\nxe2x80x9cHalogenxe2x80x9d is fluorine, chlorine, bromine or iodine, and preferably fluorine or chlorine.\nWithout limiting the present invention, in the formulae Ia, Ib, Ic, Id, Ie, If, Ig and Ih examples of groups of compounds of the invention are shown in which A2 in the formula I has specific denotations. A1, R1, R2, R3, X and n in the formulae Ia, Ib, Ic, Id, Ie, If, Ig and Ih are defined as above for the formula I, and the number k in the formula Ib is 1, 2, 3, 4 or 5, in particular 3 or 4. \nOn the benzene ring depicted in formula Ia which carries the groups C(xe2x95x90X)xe2x80x94NHxe2x80x94 and xe2x80x94NHSO2R2, four positions are present which can carry a residue R3. The compounds of formula Ia can thus carry four residues R3 which, independently of one another, can all be hydrogen or can have a meaning different from hydrogen, i.e., in the compounds of formula Ia the benzene ring depicted in formula Ia can be unsubstituted or can carry one, two, three or four identical or different substituents from the series halogen, CF3, OH, xe2x80x94Oxe2x80x94(C1-C7)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C7)-alkyl, xe2x80x94O-aryl, (C1-C2)-alkylenedioxy, NO2, xe2x80x94CN, NR7R8, xe2x80x94COxe2x80x94NR7R8, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C5)-alkyl, heterocyclyl, xe2x80x94S(O)nxe2x80x94(C1-C5)-alkyl and (C1-C5)-alkyl which can be substituted by one or more identical or different residues R4. These explanations accordingly also apply to the compounds of formulae Ib to Ih.\nThe compounds of the formula I wherein X is a nitrogen atom which via a single bond is attached to a ring carbon atom in the group A1 which ring carbon atom is directly adjacent to the carbon atom in A1 carrying the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94 so that the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94 together with the carbon atoms in A1 carrying it forms an anellated imidazole rin(g, are represented by the formula Ii. \nA2, R1, R1, R3 and nin the formula Ii are defined as above for the formula I. The ring A3 which has resulted from the group A1 by the formation of a bond to the nitrogen atom representing X and which ring comprises the two carbon atoms depicted in formula Ii carrying the nitrogen atoms of the anellated imidazole ring, is a benzene ring, a naphthalene ring or a heteroaromatic ring, where to these rings the above explanations relating to A1 correspondingly apply.\nThe present invention includes all stereoisomeric forms of the compounds of the formula I. Centers of asymmetry that are present in the compounds of formula I can all independently of one another have S configuration or R configuration. The invention includes all possible enantiomers and diastereomers and mixtures of two or more stereoisomers, for example mixtures of enantiomers and/or diastereomers, in all ratios. Thus, enantiomers are a subject of the invention in enantiomerically pure form, both as levorotatory and as dextrorotatory antipodes, in the form of racemates and in the form of mixtures of the two enantiomers in all ratios. In the case of a cis/trans isomerism the invention includes both the cis form and the trans form as well as mixtures of these forms in all ratios. The preparation of individual stereoisomers can be carried out, if desired, by separation of a mixture by customary methods, for example by chromatography or crystallization, by the use of stereochemically uniform starting materials for the synthesis or by stereoselective synthesis. Optionally a derivatization can be carried out before a separation of stereoisomers. The separation of a mixture of stereoisomers can be carried out at the stage of the compounds of the formula I or at the stage of an intermediate during the synthesis. The present invention also includes all tautomeric forms of the compounds of formula I.\nIf the compounds of the formula I contain one or more acidic or basic groups the invention also includes the corresponding physiologically or toxicologically acceptable salts, in particular the pharmaceutically utilizable salts. Thus, the compounds of the formula I which contain acidic groups can be present on these groups and can be used according to the invention, for example, as alkali metal salts, alkaline earth metal salts or as ammonium salts. Examples of such salts are sodium salts, potassium salts, calcium salts, magnesium salts or salts with ammonia or organic amines such as, for example, ethylamine, ethanolamine, triethanolamine or amino acids. Compounds of the formula I which contain one or more basic groups, i.e. groups which can be protonated, can be present and can be used according to the invention in the form of their acid addition salts with inorganic or organic acids, for example as salts with hydrogen chloride, hydrogen bromide, phosphoric acid, sulfuric acid, nitric acid, methanesulfonic acid, p-toluenesulfonic acid, naphthalenedisulfonic acids, oxalic acid, acetic acid, tartaric acid, lactic acid, salicylic acid, benzoic acid, formic acid, propionic acid, pivalic acid, diethylacetic acid, malonic acid, succinic acid, pimelic acid, fumaric acid, maleic acid, malic acid, sulfaminic acid, phenylpropionic acid, gluconic acid, ascorbic acid, isonicotinic acid, citric acid, adipic acid, etc. If the compounds of the formula I simultaneously contain acidic and basic groups in the molecule the invention also includes, in addition to the salt forms mentioned, inner salts or betaines (zwitterions). Salts can be obtained from the compounds of the formula I by customary methods which are known to the person skilled in the art, for example by combination with an organic or inorganic acid or base in a solvent or dispersant, or by anion exchange or cation exchange from other salts. The present invention also includes all salts of the compounds of the formula I which, owing to low physiological compatibility, are not directly suitable for use in pharmaceuticals but which can be used, for example, as intermediates for chemical reactions or for the preparation of physiologically acceptable salts.\nThe present invention furthermore includes all solvates of compounds of the formula I, for example hydrates or adducts with alcohols, and also derivatives of the compounds of the formula I such as, for example esters, amides, prodrugs and active metabolites.\nA1 preferably is a phenylene residue or a 5-membered or 6-membered heteroarylene residue, more preferably phenylene, where all these residues can be unsubstituted or can be substituted as indicated. If the group A1 is substituted, I. e. if it carries one or more further substituents in addition to the group R1xe2x80x94S(O)n, it is preferably substituted by one or two identical or different of the substituents indicated above. Preferably a phenylene residue representing A1 is unsubstituted, I. e. besides the groups R1xe2x80x94S(O)n and C(xe2x95x90X)xe2x80x94NH it carries four hydrogen atoms. The group R1xe2x80x94S(O)n is preferably attached to a carbon atom in A1 which is not directly adjacent to the carbon atom which carries the group C(xe2x95x90X)xe2x80x94NH. If A1 is phenylene the group R1xe2x80x94S(O)n is particularly preferably located in the meta position or in the para position, more particularly preferably in the para position, with respect to the carbon atom which carries the group C(xe2x95x90X)xe2x80x94NH.\nThe ring A2 which comprises the two carbon atoms which carry the groups R2xe2x80x94SO2xe2x80x94NH and C(xe2x95x90X)xe2x80x94NHxe2x80x94 preferably is an aromatic ring, more preferably a benzene ring or a thiophene ring, particularly preferably a benzene ring, where all these rings can be unsubstituteci or substituted by one ore more residues R3 which are different from hydrogen.\nR1 preferably is (C1-C7)-alkyl, NR5R6 or aryl, more preferably NR5R6, phenyl or 5-membered or 6-membered heteroaryl, particularly preferably NR5R6, where all these residues can be unsubstituted or substituted as indicated and where, as stated above, R1 can be NR5R6 if the number n in the group R1xe2x80x94S(O)nxe2x80x94 is 2.\nR2 preferably is aryl, more preferably phenyl or 5-membered or 6-membered heteroaryl, particularly preferably phenyl or a residue of a monocyclic 5-membered or 6-membered aromatic heterocycle which contains one or two identical or different heteroatoms from the series N, O and S such as, for example, phenyl, thienyl, pyrazolyl, imidazolyl, isoxazolyl, thiazolyl, pyridyl etc., in particular phenyl or 2-thienyl, where all these residues can be unsubstituted or substituted as indicated. Preferably an aryl residue representing R2 is substituted. If an aryl residue representing R2 is substituted it is preferably substituted by one, two or three, in particular by one or two, identical or different substituents. Substituents in an aryl residue representing R2 preferably are substituents from the series halogen, CF3, xe2x80x94Oxe2x80x94CF3, NO2, xe2x80x94CN, (C1-C4)-alkyl and xe2x80x94Oxe2x80x94(C1-C4)-alkyl, more preferably substituents from the series F, Cl, Br, CF3, xe2x80x94Oxe2x80x94CF3, NO2, xe2x80x94CN, CH3 and xe2x80x94OCH3. Especially preferably a substituted aryl residue representing R2 is substituted by Cl, for example by one or two, in particular one, chlorine atoms.\nThe rings representing A2 can be unsubstituted or substituted as indicated. When they are unsubstituted they only carry residues R3 which are hydrogen. When they are substituted they carry one or more residues R3 which are different from hydrogen. Those substituent positions which do not carry a residue R3 which is different from hydrogen, carry hydrogen atoms. If the ring A2 carries one or more residues R3 which are different from hydrogen it preferably carries one or two such residues R3, in particular one such residue R3. Residues R3 which are different from hydrogen are preferably located in positions of the ring A2 which are not directly adjacent to the groups C(xe2x95x90X)xe2x80x94NH and R2xe2x80x94SO2xe2x80x94NH. If A2 is a saturated or partially unsaturated carbocycle, residues R3 which are different from hydrogen preferably are (C1-C4)-alkyl, in particular methyl. If A2 is an aromatic ring, in particular if A2 is a benzene ring, residues R3 which are different from hydrogen preferably are (C1-C3)-alkyl, halogen, (C1-C3)-alkoxy or CF3, more preferably methyl, chlorine or methoxy. If a A2 is an aromatic ring, in particular a benzene ring, it is very particularly preferred if the ring carries one chlorine atom or two methoxy groups as substituents, I. e., if one residue R3 is present which is chlorine or if two residues R3 are present which are methoxy, and the other positions on the benzene ring carry hydrogen atoms. If A2 is a benzene ring residues R3 which are different from hydrogen are preferably located in positions 4 and/or 5 (with respect to the group C(xe2x95x90X)xe2x80x94NH in the 1-position and the group R2xe2x80x94SO2xe2x80x94NH in the 2-position).\nIf a group is substituted by one or more residues R4 it is preferably substituted by one, two or three, in particular one or two, identical or different residues R4. R4 preferably is hydroxy, (C1-C4)-alkyloxy, di-((C1-C4)-alkyl)amino or heteroaryl.\nR5 and R6 preferably are independently of one another hydrogen, (C1-C9)-alkyl, (C1-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl- or 5-membered or 6-membered aryl or together with the nitrogen atom carrying R5 and R6 form a 5-membered to 7-membered heterocycle which in addition to the nitrogen atom carrying the groups R5 and R6 can contain one further ring heteroatom from the series N, O and S and which can be substituted by one or more, for example one, two, three or four, identical or different residues from the series (C1-C3)-alkyl, hydroxy-(C1-C3)-alkyl-, 5-membered or 6-membered aryl, carbamoyl, hydroxy and oxo. It is particularly preferred if R5 and R6 together with the nitrogen atom carrying these residues form a 5-membered, 6-membered or 7-membered heterocycle which in addition to the nitrogen atom carrying the groups R5 and R6 can contain one further ring heteroatom from the series N, O and S and which can be substituted by one or more, for example one, two, three or four, identical or different residues from the series (C1-C3)-alkyl, hydroxy-(C1-C3)-alkyl-, 5-membered or 6-membered aryl, carbamoyl, hydroxy and oxo, in particular (C1-C3)-alkyl such as, for example, methyl. Preferably a heterocycle which is formed from the groups R5 and R6 together with the nitrogen atom carrying these residues is saturated.\nAn especially preferable heterocycle which is formed by R5 and R6 together with the nitrogen atom carrying these residues, is derived from morpholine, thiomorpholine, 1,1-dioxo-thiomorpholine, 1-oxo-thiomorpholine, from dialkylmorpholines such as dimethylmorpholines, from 2,6-dimethylmorpholine, cis-2,6-dimethylmorpholine, 3,5-dimethylmorpholine, cis-3,5-dimethylmorpholine, 1-(pyrimidin-2-yl)-piperazine, piperidine-4-carboxamide, 1-(2-hydroxyethyl)-piperazine, 1-methylpiperazine, 1-ethylpiperazine, from 1-arylpiperazines, from ethyl piperazine-1-carboxylate, piperidine, 2-methylpiperidine, 2,6-dimethylpiperidine, cis-2,6-dimethylpiperidine, 3,5dimethylpiperidine, cis-3,5-dimethylpiperidine, 4-hydroxypiperidine, from 4-oxopiperidine or a ketal thereof like 1,4-dioxa-8-aza-spiro[4.5]decan, from tetrahydropyridine, tetrahydropyrimidine, 1-methylhomopiperazine, thiazolidine, pyrroline, pyrrolidine, 3-hydroxypyrrolidine, 1,2,3,4-tetrahydroisoquinoline or 2,3-dihydro-1H-isoindole, where these rings are attached via the ring nitrogen atom or, in the case of piperazine derivatives, via the unsubstituted ring nitrogen atom. A more especially preferable heterocycle which is formed by R5 and R6 together with the nitrogen atom carrying these residues, is derived from morpholine, thiomorpholine, 1,1-dioxo-thiomorpholine, 1-oxo-thiomorpholine, 2,6-dimethylmorpholine, cis-2,6-dimethylmorpholine, 3,5-dimethylmorpholine, cis-3,5-dimethylmorpholine, 1-(pyrimidin-2-yl)-piperazine, piperidine-4-carboxamide, 1,2,3,4-tetrahydroisoquinoline or 2,3-dihydro-1H-isoindole, moreover preferably from morpholine, 2,6-dimethylmorpholine or cis-2,6-dimethylmorpholine, in particular from morpholine or cis-2,6-dimethylmorpholine, where these rings are attached via the ring nitrogen atom or, in the case of the piperazine derivative, via the unsubstituted ring nitrogen atom.\nR7 preferably is hydrogen, (C1-C3)-alkyl, ((C1-C4)-alkyl)2Nxe2x80x94(C1-C3)-alkyl- or (C1C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl-.\nR8 preferably is hydrogen, (C1-C3)-alkyl or acetyl.\nAryl preferably is phenyl or heteroaryl, in particular phenyl or 5-membered or 6-membered heteroaryl. Unless stated otherwise, preferred substituents on aryl residues are halogen, CF3, (C1-C3)-alkyl, cyano, nitro and (C1-C3)-alkyloxy, more preferred substituents are CF3, chlorine, methyl and methoxy.\nHeteroaryl and heteroarylene preferably are a residue of a monocyclic 5-membered or 6-membered aromatic heterocycle, in particular a residue derived from the heteroaromatics thiophene, pyrazole, thiazole, oxazole, isoxazole, pyridine, pyrimidine, pyridazine and tetrazole.\nHeterocyclyl preferably is a residue derived from a saturated heterocycle, more preferably a residue of a monocyclic 5-membered or 6-membered saturated heterocycle, in particular a residue which is derived from pyrrolidine, piperidine, from N-alkylpiperazines, from morpholine, from dialkylmorpholines, from thiomorpholine or tetrahydrofuran. In addition, the above explanations on preferred heterocycles which are formed by the residues R5 and R6 together with the nitrogen atom carrying these residues correspondingly apply to heterocyclyl residues which are attached via a ring nitrogen atom.\nIf a group S(O)n is bonded to a nitrogen atom the number n therein preferably is 1 or 2, more preferably 2. The number n in the group R1xe2x80x94S(O)n preferably is 0 or 2, particularly preferably 2.\nX preferably is O or a nitrogen atom which via a single bond is attached to a ring carbon atom in the group A1 which ring carbon atom is directly adjacent to the carbon atom in A1 carrying the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94 so that the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94 together with the carbon atoms in A1 carrying it forms an anellated imidazole ring. Particularly preferably X is O.\nPreferred compounds of the formula I are those compounds in which one or more of the residues contained therein have preferred meanings, all combinations of preferred substituent definitions being a subject of the present invention. Also with respect to all preferred compounds of the formula I the present invention includes all stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts. Groups of preferred compounds are also formed, for example, by the compounds of the formulae Ia, Ib, Ic, Id, Ie, If, Ig and Ih in which one or more residues have preferred meanings, in all their stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts.\nA group of particularly preferred compounds is formed, for example, by compounds of the formula I in which\nA1 is phenylene or heteroarylene where these residues can be unsubstituted or substituted by one or more identical or different substituents from the series halogen, (C1-C4)-alkyl, CF3, xe2x80x94Oxe2x80x94(C1-C4)-alkyl and xe2x80x94CN;\nthe ring A2 which comprises the two carbon atoms which carry the groups R2xe2x80x94SO2xe2x80x94NH and C(xe2x95x90X)xe2x80x94NHxe2x80x94 is an aromatic ring;\nR1 is (C1-C7)-alkyl which can be substituted by one or more identical or different residues R4, or is aryl, or if the number n in the group R1xe2x80x94S(O)nxe2x80x94 is 2 also is NR5R6;\nR2 is aryl;\nR3 denotes one or more identical or different residues from the series hydrogen, halogen, CF3, OH, xe2x80x94Oxe2x80x94(C1-C4)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C4)-alkyl, xe2x80x94O-aryl, NO2, xe2x80x94CN, NR7R8, xe2x80x94COxe2x80x94NR7R8, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C4)-alkyl, heterocyclyl, xe2x80x94S(O)nxe2x80x94(C1-C4)-alkyl and (C1-C4)-alkyl which can be substituted by one or more identical or different residues R4;\nR4 is fluorine, OH, xe2x80x94Oxe2x80x94(C1-C10)-alkyl, xe2x80x94Oxe2x80x94(C2-C4)-alkyl-Oxe2x80x94(C1-C7)-alkyl, xe2x80x94O-aryl, xe2x80x94CN, NR7R8, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94NHxe2x80x94(C1-C3)-alkyl, xe2x80x94COxe2x80x94N((C1-C3)-alky)2, xe2x80x94COxe2x80x94OH, xe2x80x94COxe2x80x94Oxe2x80x94(C1-C4)-alkyl, heterocyclyl or oxo;\nR5 and R6 independently of one another are hydrogen, (C1-C9)-alkyl, (C1-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl- or aryl or together with the nitrogen atom carrying R5 and R6 form a 5-membered to 7-membered heterocycle which in addition to the nitrogen atom carrying the groups R5 and R6 can contain one further ring heteroatom from the series N, O and S and which can be substituted by one or more identical or different residues from the series (C1-C3)-alkyl, hydroxy-(C1-C3)-alkyl-, aryl, carbamoyl, hydroxy and oxo;\nR7 is hydrogen, (C1-C3)-alkyl, ((C1-C4)-alkyl)2Nxe2x80x94(C1-C3)-alkyl- or (C1-C4)-alkyl-Oxe2x80x94(C1-C3)-alkyl-;\nR8 is hydrogen, (C1-C3)-alkyl or acetyl;\nxe2x80x9carylxe2x80x9d is phenyl or heteroaryl which can all be substituted by one or more identical or different substituents from the series halogen, (C1-C4)-alkyl, phenyl, CF3, NO2, xe2x80x94Oxe2x80x94(C1-C4)-alkyl, (C1-C2)-alkylenedioxy, NH2, xe2x80x94NHxe2x80x94COxe2x80x94(C1-C4)-alkyl, xe2x80x94CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94OH and xe2x80x94COxe2x80x94Oxe2x80x94(C1-C4)-alkyl;\nheteroaryl and heteroarylene are a residue of a monocyclic 5-membered or 6-membered aromatic heterocycle which contains one or more identical or different ring heteroatoms from the series N, O and S;\nheterocyclyl is a residue of a monocyclic 5-membered or 6-membered saturated heterocycle which contains one or more identical or different ring heteroatoms from the series N, O and S and which can be substituted by one or more identical or different substituents from the series fluorine, (C1-C4)-alkyl, OH, xe2x80x94Oxe2x80x94(C1-C4)-alkyl, NH2, xe2x80x94CN, xe2x80x94COxe2x80x94NH2, xe2x80x94COxe2x80x94OH and xe2x80x94COxe2x80x94Oxe2x80x94(C1-C4)-alkyl;\nn is 0, 1 or 2;\nX is oxygen;\nin all their stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts.\nA group of more particularly preferred compounds is formed, for example, by compounds of the formula I in which:\nA1 is phenylene which is unsubstituted or substituted by one or more identical or different substituents from the series halogen, (C1-C4)-alkyl, CF3, xe2x80x94Oxe2x80x94(C1-C4)-alkyl and xe2x80x94CN;\nthe ring A2 which comprises the two carbon atoms which carry the groups R2xe2x80x94SO2xe2x80x94NH and C(xe2x95x90X)xe2x80x94NHxe2x80x94 is a benzene ring;\nR1 is NR5R6;\nR2 is aryl;\nR3 denotes one or more identical or different residues from the series hydrogen, halogen, CF3, xe2x80x94Oxe2x80x94(C1-C4)-alkyl, xe2x80x94CN and (C1-C4)-alkyl;\nR5 and R 6 together with the nitrogen atom carrying R5 and R6 form a 5-membered or 6-membered saturated heterocycle which in addition to the nitrogen atom carrying the groups R 5and R6 can contain one further ring heteroatom from the series N, O and S and which can be substituted by one or more identical or different residues from the series (C1-C3)-alkyl, hydroxy-(C1-C3)-alkyl-, aryl, carbamoyl, hydroxy and oxo;\naryl is phenyl or 5-membered or 6-membered heteroaryl containing one or more identical or different ring heteroatoms from the series N, O and S which residues can all be substituted by one or more identical or different substituents from the series halogen, (C1-C4)-alkyl, CF3, NO2, xe2x80x94Oxe2x80x94(C1-C4)-alkyl, xe2x80x94NHxe2x80x94COxe2x80x94(C1-C4)-alkyl and xe2x80x94CN;\nn is 2;\nX is oxygen;\nin all their stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts.\nA group of especially preferred compounds is formed, for example, by compounds of the formula I in which:\nA1 is an unsubstituted divalent phenylene residue;\nthe ring A2 which comprises the two carbon atoms which carry the groups R2xe2x80x94SO2xe2x80x94NH and C(xe2x95x90X)xe2x80x94NHxe2x80x94 is a benzene ring;\nR1 is NR 5R6;\nR2 is aryl;\nR3 denotes one or more identical or different residues from the series hydrogen, halogen, xe2x80x94Oxe2x80x94(C1-C4)-alkyl and (C1-C4)-alkyl;\nR5 and R6 together with the nitrogen atom carrying R5 and R6 form a saturated 6-membered heterocycle which in addition to the nitrogen atom carrying the groups R5 and R6 can contain one further ring heteroatom from the series N, O and S and which can be substituted by one or more identical or different residues from the series (C1-C3)-alkyl, aryl, oxo and carbamoyl;\naryl is phenyl or 5-membered or 6-membered heteroaryl containing one or more identical or different ring heteroatoms from the series N, O and S which residues can all be substituted by one or more identical or different substituents from the series halogen, (C1-C4)-alkyl, CF3 and xe2x80x94Oxe2x80x94(C1-C4)-alkyl;\nn is 2;\nX is oxygen;\nin all their stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts.\nA group of even more preferred compounds is formed, for example, by compounds of the formula I in which:\nA1 is an unsubstituted divalent 1,4-phenylene residue;\nthe ring A2 which comprises the two carbon atoms which carry the groups R2xe2x80x94SO2xe2x80x94NH and C(xe2x95x90X)xe2x80x94NHxe2x80x94, together with the residues R3, is a benzene ring which carries one or two substituents from the series chlorine and methoxy;\nR1 is NR5R6;\nR2 is phenyl or thienyl which residues are all substituted by one or two chlorine atoms;\nR5 and R6 together with the nitrogen atom carrying R5 and R6 form a saturated 6-membered heterocycle which in addition to the nitrogen atom carrying the groups R5 and R6 can contain one further ring heteroatom from the series O and S and which is unsubstituted or substituted by one or two methyl residues;\nn is 2;\nX is oxygen;\nin all their stereoisomeric forms and mixtures thereof in all ratios, and their physiologically acceptable salts.\nThe presents invention also relates to processes for the preparation of the compounds of the formula I which are described in the following and by which the compounds of the invention are obtainable. According to Scheme 1 compounds of the invention can be obtained, for example, by first reacting an aminocarboxylic acid of the formula II with a sulfonyl chloride of the formula R2xe2x80x94SO2xe2x80x94Cl or a sulfonic acid anhydride in the presence of a base in a solvent like water, pyridine or an ether. Suitable bases are inorganic bases like sodium carbonate or organic bases like, for example, pyridine or triethylamine. The sulfonylaminocarboxylic acid of the formula III that is obtained can then be activated, for example by reaction with a chlorinating agent like phosphorus pentachloride, phosphorus oxychloride or thionyl chloride in an inert solvent, to give an acid chloride of the formula IV which is then reacted with an arylamine.\nThe activation of the carboxylic acid group in the compounds of the formula III can, however, also be carried out in a different manner, for example by one of the numerous methods for the formation of amide bonds in peptide chemistry which are well-known to the skilled person, for example by conversion into a mixed anhydride or an activated ester or by using a carbodiimide like dicyclohexylcarbodimide.\nThe reaction of the activated sulfonylaminocarboxylic acid with an arylamine is favorably carried out in an inert solvent such as, for example, pyridine, tetrahydrofuran or toluene in the absence or in the presence of an inert auxiliary base like, for example, a tertiary amine or pyridine. If the arylamine that is employed in the reaction with the activated carboxylic acid already contains the desired substituent R1xe2x80x94S(O)n then the reaction directly provides the final compound of the formula I. Compounds of the formula I in which the number n in the group R1xe2x80x94S(O)n is 1 or 2 can also be obtained by reacting an activated carboxylic acid with a mercapto substituted arylamine of the formula R1xe2x80x94Sxe2x80x94A1xe2x80x94NH2 and then oxidizing the mercapto group in the compound of the formula V under standard conditions, for example with a peroxide like hydrogen peroxide or a peracid like 3-chloroperbenzoic acid or monoperoxyphthalic acid in a solvent like, for example, methylene chloride or acetone. The activated carboxylic acids can also first be reacted with arylamines of the formula A1xe2x80x94NH2. The resulting reaction product of the formula VI can then be chlorosulfonated under standard conditions and the chlorosulfonyl group can then be converted under standard conditions into the group R1xe2x80x94SO2, for example by reaction with suitable amines in substance or in a solvent like N-methylpyrrolidone, dimethylformamide, toluene or an ether, optionally in the presence of an auxiliary base. In a similar manner the activated carboxylic acids can be reacted with fluorosulfonylarylamines of the formula Fxe2x80x94SO2xe2x80x94A1xe2x80x94NH2 and the fluorosulfonyl intermediates of the formula VII that are obtained can be converted under standard conditions into the compounds of the formula I according to the invention. \nThe compounds of the formula I according to the invention can furthermore be obtained by reacting the activated sulfonylaminocarboxylic acids, for example the acid chlorides of the formula IV shown in Scheme 1, with a mercaptoarylamine of the formula H2Nxe2x80x94A1xe2x80x94SH which is unsubstituted on the sulfur atom. In a nucleophilic substitution reaction the product of the formula VIII that is obtained can subsequently be alkylated or arylated on the sulfur atom with an alkyl halogenide or an aryl halogenide or another reactive compound under standard conditions and, if desired, be oxidized on the sulfur to give the sulfoxide or the sulfone as explained above with respect to the compounds of the formula V (see Scheme 2). \nCompounds of the formula I can also be obtained, for example, by first activating a suitably substituted nitrocarboxylic acid of the formula IX, for example by converting it into the respective acid chloride of the formula X or by another procedure, and then reacting it with a substituted arylamine of the formula R1xe2x80x94S(O)nxe2x80x94A1xe2x80x94NH2 analogously to the procedures described above (see Scheme 3). Here, too, as arylamines fluorosulfonylarylamines of the formula Fxe2x80x94SO2xe2x80x94A1xe2x80x94NH2 can be employed, and in the N-(fluorosulfonylaryl)-carboxamides of the formula XI that are obtained the fluorosulfonyl group can be converted under standard conditions into a group R1xe2x80x94SO2 according to the invention, for example with an amine of the formula HNR5R6. \nBefore the nitro group is reduced to an amino group, in the nitro intermediates of the formula XII, use can be made of the activating effect of the nitro group on the ring A2 and a suitable residue R3, for example a halogen atom, can be replaced with a different residue R3 by reaction with a nucleophile, for example an amine. The reduction of the nitro group to give an amino group can, for example, be carried out by catalytic hydrogenation in the presence of a noble metal catalyst or, preferably, in the presence of Raney nickel in a solvent like ethanol, glacial acetic acid or an ethanolic solution of hydrogen chloride, or it can be carried out by reduction with a base metal like zinc, tin or iron in the presence of an acid. The reduction can also be carried out, for example, with tin(II) chloride or by reaction with sodium dithionite, favorably in a mixture of methanol, tetrahydrofuran and water as solvent. The sulfonylation of the amino group in the reduction product of the formula XIII with an activated sulfonic acid derivative can be carried out analogously to the reactions described above, for example with a sulfonic acid chloride in pyridine, and finally gives the compound of the formula I.\nThe compounds of the formula I in which X is a nitrogen atom which via a single bond is attached to a ring carbon atom in the group A1 which ring carbon atom is directly adjacent to the carbon atom in A1 carrying the group xe2x80x94NHxe2x80x94C(xe2x95x90X)xe2x80x94, I. e. the benzimidazole derivatives of the formula Ii, can for example be obtained by reacting an activated sulfonylaminocarboxylic acid derivative obtained as described above according to Scheme 1, for example a carboxylic acid chloride of the formula IV, (or also, in analogy to Scheme 3, a nitrocarboxylic acid derivative) with a 1,2-diaminoarene in the presence of a dehydrating agent such as, for example, thionyl chloride or phosphorus pentachloride (see Scheme 4). The reaction is usually carried out in an inert solvent, for example in a hydrocarbon like toluene or xylene. The 1,2-diaminoarene can already contain the final group R1xe2x80x94S(O)n or a precursor group thereof, for example the group R1xe2x80x94S. Subsequent steps, for example reactions on the sulfur atom, can then be carried out as explained above. Just so, unsubstituted 1,2-diaminoarenes can be employed and the resulting products of the formula XIV can be chlorosulfonated, for example with chlorosulfuric acid, and the sulfonyl chlorides that are obtained can be converted into the final compounds containing the group R1xe2x80x94SO2, for example by reaction with a suitable amine. \nAll reactions for the synthesis of the compounds of the formula I are per se well-known to the skilled person and can be carried out under standard conditions according to or analogously to procedures described in the literature, for example in Houben-Weyl, Methoden der Organischen Chemie (Methods of Organic Chemistry), Thieme-Verlag, Stuttgart, or Organic Reactions, John Wiley and Sons, New York. Depending on the circumstances of the individual case, in order to avoid side reactions during the synthesis of a compound of the formula I it can be necessary or advantageous to temporarily block functional groups by introducing protective groups and to deprotect them in a later stage of the synthesis, or introduce functional groups in the form of precursor groups which in a later reaction step are converted into the desired functional groups. Such synthetic strategies and protective groups and precursor groups which are suitable in an individual case are known to the skilled person. If desired the compounds of the formula I can be purified by customary purification procedures, for example by recrystallization or chromatography. The starting compounds for the preparation of the compounds of the formula I are commercially available or can be prepared according to or analogously to literature procedures.\nThe compounds of the formula I according to the invention effect an increase of the cGMP concentration via the activation of the soluble guanylate cyclase (sGC), and they are therefore useful agents for the therapy and prophylaxis of disorders which are associated with a low or decreased cGMP level or which are caused thereby, or for whose therapy or prophylaxis an increase of the present cGMP level is desired. The activation of the sGC by the compounds of the formula I can be examined, for example, in the activity assay described below.\nDisorders and pathological conditions which are associated with a low cGMP level or in which an increase of the cGMP level is desired and for whose therapy and prophylaxis it is possible to use compounds of the formula I are, for example, cardiovascular diseases, such as endothelial dysfunction, diastolic dysfunction, atherosclerosis, hypertension, stable and unstable angina pectoris, thromboses, restenoses, myocardial infarction, strokes, cardiac insufficiency or pulmonary hypertonia, or, for example, erectile dysfunction, asthma bronchiale, chronic kidney insufficiency and diabetes. Compounds of the formula I can additionally be used in the therapy of cirrhosis of the liver and also for improving a restricted memory performance or ability to learn.\nThe compounds of the formula I and their physiologically acceptable salts can be administered to animals, preferably to mammals, and in particular to humans, as pharmaceuticals by themselves, in mixtures with one another or in the form of pharmaceutical preparations. A subject of the present invention therefore also are the compounds of the formula I and their physiologically acceptable salts for use as pharmaceuticals, their use for activating soluble guanylate cyclase, for normalizing a disturbed cGMP balance and in particular their use in the therapy and prophylaxis of the abovementioned syndromes as well as their use for preparing medicaments for these purposes.\nFurthermore, a subject of the present invention are pharmaceutical preparations (or pharmaceutical compositions) which comprise as active component an effective dose of at least one compound of the formula I and/or a physiologically acceptable salt thereof and a customary pharmaceutically acceptable carrier, I. e. one or more pharmaceutically acceptable carrier substances and/or additives. A subject of the present invention is also those compounds of the formula I which were already known per se and which are excluded by disclaimer from the above-defined compounds of the formula I which are per se a subject of the present invention, and their physiologically acceptable salt as activators of soluble guanylate cyclase.\nAll statements above and below relating to the pharmacological effects and the uses of the compounds of the formula I thus also apply to the compound of the formula I wherein simultaneously the ring A2 which comprises the carbon atoms which carry the groups C(xe2x95x90X)xe2x80x94NHxe2x80x94 and NHxe2x80x94SO2R2 is a benzene ring which is substituted in positions 3 and 5 by chlorine, R2 is methyl, X is oxygen and R1xe2x80x94S(O)nxe2x80x94A1xe2x80x94 is a 5-chloro-2-(4-chlorophenylmercapto)-phenyl residue, and its physiologically acceptable salts.\nThus, a subject of the invention are, for example, said compound and its physiologically acceptable salts for use as a pharmaceutical, pharmaceutical preparations which comprise as active component an effective dose of said compound and/or a physiologically acceptable salt thereof and a customary pharmaceutically acceptable carrier, and the uses of said compound and/or a physiologically acceptable salt thereof in the therapy or prophylaxis of the abovementioned syndromes as well as their use for preparing medicaments for these purposes.\nThe pharmaceuticals according to the invention can be administered orally, for example in the form of pills, tablets, lacquered tablets, sugar-coated tablets, granules, hard and soft gelatin capsules, aqueous, alcoholic or oily solutions, syrups, emulsions or suspensions, or rectally, for example in the form of suppositories. Administration can also be carried out parenterally, for example subcutaneously, intramuscularly or intravenously in the form of solutions for injection or infusion. Other suitable administration forms are, for example, percutaneous or topical administration, for example in the form of ointments, tinctures, sprays or transdermal therapeutic systems, or the inhalative administration in the form of nasal sprays or aerosol mixtures, or, for example, microcapsules, implants or rods. The preferred administration form depends, for example, on the disease to be treated and on its severity.\nThe amount of active compound of the formula I and/or its physiologically acceptable salts in the pharmaceutical preparations normally is from 0.2 to 500 mg, preferably from 1 to 200 mg, per dose, but depending on the type of the pharmaceutical preparation it can also be higher. The pharmaceutical preparations usually comprise 0.5 to 90 percent by weight of the compounds of the formula I and/or their physiologically acceptable salts. The preparation of the pharmaceutical preparations can be carried out in a manner known per se. For this purpose, one or more compounds of the formula I and/or their physiologically acceptable salts, together with one or more solid or liquid pharmaceutical carrier substances and/or additives (or auxiliary substances) and, if desired, in combination with other pharmaceutically active compounds having therapeutic or prophylactic action, are brought into a suitable administration form or dosage form which can then be used as a pharmaceutical in human or veterinary medicine.\nFor the production of pills, tablets, sugar-coated tablets and hard gelatin capsules it is possible to use, for example, lactose, starch, for example maize starch, or starch derivatives, talc, stearic acid or its salts, etc. Carriers for soft gelatin capsules and suppositories are, for example, fats, waxes, semisolid and liquid polyols, natural or hardened oils, etc. Suitable carriers for the preparation of solutions, for example of solutions for injection, or of emulsions or syrups are, for example, water, physiologically sodium chloride solution, alcohols such as ethanol, glycerol, polyols, sucrose, invert sugar, glucose, mannitol, vegetable oils, etc. It is also possible to lyophilize the compounds of the formula I and their physiologically acceptable salts and to use the resulting lyophilisates, for example, for preparing preparations for injection or infusion. Suitable carriers for microcapsules, implants or rods are, for example, copolymers of glycolic acid and lactic acid.\nBesides the active compounds and carriers, the pharmaceutical preparations can also contain customary additives, for example fillers, disintegrants, binders, lubricants, wetting agents, stabilizers, emulsifiers, dispersants, preservatives, sweeteners, colorants, flavorings, aromatizers, thickeners, diluents, buffer substances, solvents, solubilizers, agents for achieving a depot effect, salts for altering the osmotic pressure, coating agents or antioxidants.\nThe dosage of the active compound of the formula I to be administered and/or of a physiologically acceptable salt thereof depends on the individual case and is, as is customary, to be adapted to the individual circumstances to achieve an optimum effect. Thus, it depends on the nature and the severity of the disorder to be treated, and also on the sex, age, weight and individual responsiveness of the human or animal to be treated, on the efficacy and duration of action of the compounds used, on whether the therapy is acute or chronic or prophylactic, or on whether other active compounds are administered in addition to compounds of the formula I. In general, a daily dose of approximately 0.01 to 100 mg/kg, preferably 0.1 to 10 mg/kg, in particular 0.3 to 5 mg/kg (in each case mg per kg of bodyweight) is appropriate for administration to an adult weighing approximately 75 kg in order to obtain the desired results. The daily dose can be administered in a single dose or, in particular when larger amounts are administered, be divided into several, for example two, three or four individual doses. In some cases, depending on the individual response, it may be necessary to deviate upwards or downwards from the given daily dose.\nThe compounds of the formula I activate the soluble guanylate cyclase, mainly by binding in the heme binding pocket of the enzyme. On account of this property, apart from use as pharmaceutically active compounds in human medicine and veterinary medicine, they can also be employed as a scientific tool or as aids for biochemical investigations in which such an effect on guanylate cyclase is intended, and also for diagnostic purposes, for example in the in vitro diagnosis of cell samples or tissue samples. The compounds of the formula I and salts thereof can furthermore be employed, as already mentioned above, as intermediates for the preparation of other pharmaceutically active compounds.\nThe following examples of compounds of the formula I and of intermediates for their preparation illustrate the invention without limiting it."}
{"text": "A large number of electric connections, particularly electrical plug connections, are known which serve to transmit electric currents, voltages and/or signals with a largest possible bandwidth. Particularly in the automotive industry, such connections must safeguard a faultless transmission of electric power, signals and/or data in thermally charged, polluted, moist or chemically aggressive surroundings.\nDue to a wide range of applications for such connections, a large number of specifically configured electric plug contacts are known, particularly crimp-contacts. In the field of electrical power contacting for the automotive industry, aside from a crimp-contact, only circular high-voltage or high-current contacts are known which could easily be stamped out of milled metal strips. In a rectangular high-voltage or high-current contact, an electric contact is provided by many filigree contact lamellas, wherein all contact lamellas have the same design and are bound to a contact cage at both longitudinal end portions. Due to the position of the contact lamellas in the contact, an amperage varies per contact lamella; a balanced current distribution is not possible with such a contact. Furthermore, the many filigree contact lamellas lead to a non-robust, damageable contact.\nA known contact comprises identical contact lamellas wherein some contact lamellas are located more closely to a conductor-crimping section of the contact than several other contact lamellas. When using the contact, because the current always takes the path of least resistance, this leads to the problem that the contact lamellas which are located more closely to the conductor-crimping section carry more electric current than those which are located further away from the conductor-crimping section. The contact lamella located closest to the conductor-crimping section carries the most current and the one furthest away from the conductor-crimping section only carries a very small amount or hardly any current."}
{"text": "The inventors have discovered, during the course of their histological studies on antitumor agents, the novel antitumor antibiotic BMG 162-aF2, later named spergualin, in a broth for culture of Bacillus laterosporus BMG 162-aF2 (Deposit No. 5230 at the Fermentation Research Institute, Agency of Industrial Science and Technology, Ministry of International Trade and Industry of Japan), a strain belonging to the genus Bacillus (Journal of Antibiotics, Vol. 34, 1619-1622, 1981), The chemical structure of spergualin is represented by the following formula ##STR3## in which the configuration at the 15-position is S-configuration (sinister), but that at the 11-position has not been determined (Journal of Antibiotics, vol. 34, 1622, 1981).\nThe compound of this structural formula could also be synthesized by the condensation of an acid amide and glyoxylylspermidine, and the resulting epimer was divided into (-)-spergualin that naturally occurs and (+)-spergualin that does not (Journal of Antibiotics, Vol. 34, 1625, 1981)."}
{"text": "Embodiments of the present disclosure relate to three-dimensional (3D) memory devices and fabrication methods thereof.\nPlanar memory cells are scaled to smaller sizes by improving process technology, circuit design, programming algorithm, and fabrication process. However, as feature sizes of the memory cells approach a lower limit, planar process and fabrication techniques become challenging and costly. As a result, memory density for planar memory cells approaches an upper limit.\nA 3D memory architecture can address the density limitation in planar memory cells. The 3D memory architecture includes a memory array and peripheral devices for controlling signals to and from the memory array."}
{"text": "(5.1) Field of the Invention\nThe present invention relates generally to computer programs and, more specifically to a method and apparatus for processing escrow transactions and a method of doing business in escrow transactions using computers and computer network technology. More particularly, the present invention relates to a computerized method and apparatus for escrow transactions and a method of doing business in escrow transaction management, preferably using the Internet.\n(5.2) Description of Related Art\n(5.2.1) Glossary\nThe following terms and definitions are provided for the convenience of the reader; no limitation on the scope of the invention is intended by this GLOSSARY nor should any be implied therefrom.\n“Client-Server”: A model of interaction in a distributed computer network system in which a program at one site sends a request to another site and then waits for a response. The requesting program is called the “client,” and the program which responds to the request is called the “server.” In the context of the World Wide Web (“www” or “web” defined hereinafter), the client is a “browser;” i.e., a program which runs on a computer of an end-user. A program and network computer which responds to a browser request by serving web pages and the like, is referred to as a “server.” Specialized servers, such as dedicated electronic mail (defined hereinafter) servers are also known in the art.\n“Electronic Mail” (“e-mail”): The process and software for sending and receiving of textual information and attachments thereto between end-users over a distributed computer network such as the internet; internet access providers often include e-mail service to its customers as part of the access software that allows the end user to dial into the internet.\n“Hypertext Mark-up Language” (“HTML”): A standard coding convention and set of codes for attaching presentation and linking attributes to informational content within documents; the primary standard used for generating web documents. During a document authoring stage, the HTML codes (referred to as “tags”) are embedded within the informational content of the document; when the document is subsequently transferred from a server to a client, the codes are interpreted by the browser and used to parse and display the document. In specifying how the browser is to display the document, HTML tags can be used to create hyperlinks to other web documents.\n“HyperText Transport Protocol” (“HTTP”): The standard Internet (infra) client-server protocol used for the exchange of information such as HTML documents and client requests for such documents between a browser and the server. HTTP—or HML or the like—includes a number of different types of messages which can be sent from the client to the server to request different types of server actions. For example, a “get” message which has the format GET<URL> (defined hereinafter) causes the server to return the document or file located at the specified URL.\n“Internet”: A generic term for a collection of distributed, interconnected networks (ARPANET, DARPANET, World Wide Web, or the like) that are linked together by a set of industry standard protocols and the like to form a global, distributed network.\nWeb “Site”: A computer system that serves informational content over a network using standard protocols of the web. Typically, a web site corresponds to a particular internet domain name, such as WIDGET.COM”, and includes the content associated with a particular organization such as Widget Company. The term is generally intended to encompass both (1) the hardware/software server components that serve the informational content over the network, and (2) the “back-end” hardware-software components, including any non-standard or specialized components that interact with the server component to perform service for web site users.\n“World Wide Web” (WWW or simply “web”): Refers generally to both (1) a distributed collection of interlinked, user-viewable hypertext documents (“web documents” or “web pages”) that are accessible via the internet, and (2) the client and server software components which provide user access to such documents using standardized internet protocols. Currently, the primary standard protocol for allowing applications to locate and acquire web documents is HTTP, and the web pages are encoded using HTML. However, the terms “web” and WWW or “world wide web” as used herein are intended to encompass future mark-up languages and transport protocols which may be used.\n(5.2.2.) Background of the Invention\nIn the main, an escrow is a third party to a transaction between a seller and a buyer (the buyer and seller are collectively referred to as the “principals”). A deed, a bond, money, a piece of property, or other valuable, is delivered to the third party—hereinafter referred to as the “escrow officer”—to be delivered by the officer to the grantee only upon the fulfillment of all the conditions precedent which are imposed upon the principals and are usually of the express terms and conditions of a purchase and sale agreement (also referred to more simply as the “sales contract” hereinafter). In essence, the property is placed in trust in an escrow account.\nEscrow transactions can be for both personal and real property. Some states (e.g. California) provide for the use of licensed escrow agents when a sale of real estate is being transacted. For the purpose of describing the present invention, an exemplary escrow transaction for real estate is considered; no limitation on the scope of the invention is intended by the inventor nor should any be implied therefrom.\nWhile computers are important tools for the buyers and sellers of a valuable property (jointly referred to as the “principals”), in the main, the use of actual purchase agreements, escrow services and associated contracts, such as loan agreements, are handled manually since the true identity of the principals and their agreement by signatur on the assorted contracts and documents involved is a critical, personal factor. Notarization on the more important documents is often required. The coordination of the various entities involved with the transfer of a real estate property—namely the principals, agents, real estate brokers, attorneys, lending institutions such as banks or other mortgage related institutions, insurance companies, government and quasi-governmental entities (such as county clerk offices, local real estate boards, state deed recording and tax departments, and the like), vendors (e.g., title companies, property inspectors, property assessors, home owners associations, and the like), and any and all other entities involved in the transaction to be escrowed, is assigned to a licensed escrow agent. In turn, the escrow agent usually must involve the aid of administrative assistants to track all of the required paperwork and transfer of funds necessary to complete the transaction. As can be imagined, the number of such real estate transactions in a single California county on a daily basis is enormous. Moreover, commercial realty transactions involve even more complex service involving even more service providers.\nWith the advent of the so-called computer age, lending institution have devised methods and apparatus for processing mortgages and other loans. Mortgages and other transactions involving banks and credit-loan associations have been the subject of computerization for many years. See, e.g., U.S. Pat. Nos.\nU.S. Pat. No. 4,876,648, SYSTEM AND METHOD FOR IMPLEMENTING AND ADMINISTERING A MORTGAGE PLAN;\nU.S. Pat. No. 5,924,084, NEGOTIATED MATCHING SYSTEM;\nU.S. Pat. No. 5,930,776, LENDER DIRECT CREDIT EVALUATION AND LOAN PROCESSING SYSTEM;\nU.S. Pat. No. 5,940,811, CLOSED LOOP FINANCIAL TRANSACTION METHOD AND APPARATUS;\nU.S. Pat. No. 5,983,206, COMPUTER SYSTEM AND COMPUTER-IMPLEMENTED PROCESS FOR IMPLEMENTING A MORTGAGE PARTNERSHIP;\nU.S. Pat. No. 5,966,699, SYSTEM AND METHOD FOR CONDUCTING LOAN AUCTION OVER COMPUTER NETWORK;\nU.S. Pat. No. 5,991,745, REVERSE MORTGAGE LOAN CALCULATION SYSTEM AND PROCESS;\nU.S. Pat. No. 5,995,947, INTERACTIVE MORTGAGE AND LOAN INFORMATION AND REAL-TIME TRADING SYSTEM;\nU.S. Pat. No. 6,012,047, REVERSE MORTGAGE PROCESSING SYSTEM;\nU.S. Pat. No. 6,016,482, ENHANCED COLLATERALIZED FUNDING PROCESSOR; and\nU.S. Pat. No. 6,029,149, LENDER DIRECT CREDIT EVALUATION AND LOAN PROCESSING SYSTEM.\nBuyers are often “pre-qualified” by a lending institution for a particular loan amount even before seeking the right property. As a result, one of the problems of the marketplace is that the cause of delays experienced in valuable property exchange escrow transactions have shifted from the lending institution to the escrow agent. In some states, e.g., California, real estate transactions are regularly conducted via an escrow.\nWith the establishment of the internet, the world wide web, and the like, the emergence of e-commerce has generated a need for a computerized method and apparatus for processing escrow transactions and a method of managing escrows using such tools."}
{"text": "1. Field of the Invention\nThe present invention relates to an oscillating angular speed sensor.\n2. Description of Related Art\nJP-A-2002-500961 discloses a micro mechanical component produced by connecting a substrate wafer and a cap wafer. The substrate wafer has a surface structure to detect a physical quantity. The cap wafer has a recess opposing to the surface structure of the substrate wafer, such that the surface structure can be prevented from having a contact from outside. The cap wafer is bonded to the substrate wafer, and the surface structure is sealed in a space between the wafers.\nThe space between the wafers is a single closed space, such that the surface structure is arranged in the single space. When the surface structure operates as an oscillating angular speed sensor, the sensor includes a detector to detect an angular speed, and a driving portion to drive the detector in a single direction. In this case, the detector and the driving portion are arranged in the same single space.\nHowever, in a case that gas is filled in the space, when the driving portion drives the detector, oscillation amplitude of the detector may be affected by a damping of the gas. Therefore, it is difficult for the driving portion to drive the detector in the single direction, such that a detection accuracy of the sensor may become worse. Thus, the space may be required to be vacuumed for the driving portion, so as to secure the oscillation amplitude of the detector. However, when the space is vacuumed, the gas cannot have an effect of the damping. Therefore, electrodes of the detector may have damage due to collision or sticking, when an impact is applied to the sensor.\nThat is, a space suitable for the detector, and a space suitable for the driving portion are different from each other."}
{"text": "Currently, Chinese herbal injections for treatment of cardio-cerebral vascular diseases and fundus diseases include Ixeris Sonchifolia Hance injection, which for clinical application is a brown-yellow transparent liquid derived from Ixeris Sonchifolia Hance of Compositae and extracted as intravenous injection starting from full plants. The main pharmaceutical components of the intravenous injection are flavone and adenosine. The coexistence of phyto-flavone and -adenosine in Ixeris Sonchifolia Hance has shown a remarkable complementary effect for treatment of cardio-cerebral diseases, which has already been demonstrated by pharmacodynamics study and clinical application. The pharmacological effects of Ixeris Sonchifolia Hance injection consist in its efficacies of: 1. increasing coronary artery flow, lowering cardiovascular resistance, resisting myocardial infarction, enhancing collateral circulations, reducing myocardial oxygen consumption and improving cardio-microcirculations, which are useful for the treatment of coronary heart diseases, angina pectoris, chest distress and breath shortness, and myocardial infarction; 2. reducing platelet conglomeration, increasing the activity of fibrinolysin, inhibiting thrombosis, decreasing the viscosity of blood plasma and serum, increasing the electrophoresis velocity of erythrocytes, lowering cerebro-vascular resistance, increasing cerebral blood capacity and promoting restoration of neural function, which are useful in the treatment of cerebral infarction (cerebral thrombosis); 3. improving microcirculation disorders caused by bacteria, systemic microcirculation disorders caused by polymer dextran, and fundus microcirculation, and dilating fundus artery, which are useful in the treatment of fundus diseases, such as central retinitis, optic atrophy and retinitis pigmentosa. However, Ixeris Sonchifolia Hance injections exhibit poor stability and can only stand short-term storage. The contents of flavone and adenosine have been determined by High Performance Liquid chromatography (HPLC) and derivatization method, which show that adenosine content reduces from 15.0 μg/ml prior to finishing of the aqueous injection to 6.7 μg/ml thereafter, while for flavone, from 0.25 mg/ml to 0.169 mg/ml, with a lose of 55.3% and 32.4% of the prescribed, respectively. Six months later, adenosine content reduces to 6.5 μg/ml, while flavone content is 0.133 mg/ml. This can also be verified by subjecting Ixeris Sonchifolia Hance aqueous injection wherein flavone and adenosine contents are respectively 5.07 mg/ml and 24.37 μg/ml, to ten-day accelerated stress test carried out at 80° C. in an oven, followed by content determination of both, in doing so, merely 3.16 mg/ml and 12.18 μg/ml are left for flavone and adenosine, respectively (Tab. 1). In view of the considerable lose of both active ingredients of flavone and adenosine during preparation and storage, the therapeutic effect of the injection is seriously impacted. Moreover, Ixeris Sonchifolia Hance injection exhibits a distinct change in color of solution. The absorbance measured at 400 nm using ultraviolet spectrophotometer during preparation of the aqueous injection from stock extract, ranges from 0.338 prior to sterilization to 0.423 thereafter, and further increases to 0.443 after 6 months with darkened color. It should be further noted that HPLC spectra of the injection exhibit poor similarity, that is to say, there exist major differences among HPLC spectra of different batches of aqueous injections yet formulated from the same stock extract. The stability behavior of same kind of commercial aqueous injections by random sampling (Batch No. 20000303, 20000503, 20010120) shows that: flavone contents are 0.27 mg/ml, 0.51 mg/ml and 0.30 mg/ml, while for adenosine, the contents are 1.11 μg/ml, 0.147 μg/ml and 0.00 μg/ml, respectively\nIn a conclusion, Ixeris Sonchifolia Hance aqueous injection is susceptible to various factors such as preparation, storage, and thus is difficult to be controlled in quality, which greatly expenses the therapeutic effect of Ixeris Sonchifolia Hance injection."}
{"text": "Content sharing services serve content, such as pictures, videos, text, or combinations thereof, to visitors who access the content sharing service. The content may be sourced from a publisher, automatically generated, or uploaded by one of the visitors to the content sharing service. The content may be represented as digitally encoded information. The content sharing service may store the content, or link to other services and subsequently retrieve the content prior to serving the content to the visitors.\nThe visitor may employ various techniques to access the content. For example, the visitor may access the content sharing service through a browser. In another example, the visitor may access the content sharing service via an application installed on a mobile device. The visitor may use any sort of device enabled to interact with the content sharing service, such as a personal computer, mobile phone, or an Internet enabled television, for example.\nThe content sharing service may rank the content based on various heuristics. For example, the content sharing service may employ a ranking based on a count associated with the accessing of a specific content item. The content sharing service may relate the count with the popularity of a specific content item. Thus, a content item that has been accessed one million times may be deemed more popular than a content item that has been accessed one thousand times.\nThe content sharing service may serve a list of the most popular content items based on the above-described ranking. The list may strategically be embedded with content already being served by the content sharing service.\nFor example, the list may be presented to a visitor to the content sharing service in response to the visitor indicating a desire to terminate access to the content sharing service (log-off). Due to the fact that the list contains content items that have proven to be successful at captivating other users' interests, the list may be successful at convincing the user to not log-off the content sharing service. In another example, the list may be presented in response to any sort of Information identifying the visitor not being available, such as a start page to the content sharing service.\nGenerally, the list serves to incentivize the content sharing service's user base to maintain access to the content sharing service. By retaining users to view additional content, the content sharing service may monetize the serving of shared content to the users.\nThe shared content may be served before, during and after the presentation of the content. The shared content may be associated with meta information, and when the shared content is clicked-through by a user, the user may be redirected to additional content associated with the shared content. The shared content may provide Information associated with a product or service related to the content."}
{"text": "Glycolic acid (HOCH2COOH; CAS Registry Number is 79-14-1) is the simplest member of the α-hydroxy acid family of carboxylic acids. Its properties make it ideal for a broad spectrum of consumer and industrial applications, including use in water well rehabilitation, the leather industry, the oil and gas industry, the laundry and textile industry, as a monomer in the preparation of polyglycolic acid (PGA), and as a component in personal care products. Glycolic acid also is a principle ingredient for cleaners in a variety of industries (dairy and food processing equipment cleaners, household and institutional cleaners, industrial cleaners [for transportation equipment, masonry, printed circuit boards, stainless steel boiler and process equipment, cooling tower/heat exchangers], and metals processing [for metal pickling, copper brightening, etching, electroplating, electropolishing]). It has also been reported that polyglycolic acid is useful as a gas barrier material (i.e., exhibits high oxygen barrier characteristics) for packing foods and carbonated drinks (WO 2005/106005 A1). However, traditional chemical synthesis of glycolic acid produces a significant amount of impurities that must be removed prior to use. New technology to commercially produce glycolic acid, especially one that produces glycolic acid in high purity and at low cost, would be eagerly received by industry.\nMicrobial enzyme catalysts can hydrolyze a nitrile (e.g., glycolonitrile) directly to the corresponding carboxylic acids (e.g., glycolic acid) using a nitrilase (EC 3.5.5.7), where there is no intermediate production of the corresponding amide (Equation 1), or by a combination of nitrile hydratase (EC 4.2.1.84) and amidase (EC 3.5.1.4) enzymes, where a nitrile hydratase (NHase) initially converts a nitrile to an amide, and then the amide is subsequently converted by the amidase to the corresponding carboxylic acid (Equation 2):\n\nIt has been demonstrated that the enzyme catalyst specific activity measured as micromoles glycolic acid produced per minute per g dry cell weight of catalyst, decreases by from 35% to 50% of initial activity after a single use in consecutive batch reactions with biocatalyst recycle. This represents a significant loss in specific activity of the enzyme catalyst in batch reactions, and, in turn, overall enzyme activity and productivity. For a commercially feasible enzymatic process for producing glycolic acid, this loss in enzyme activity needs to be addressed.\nOne aspect of the loss of enzyme activity may be attributable to impurities and other components present when reacting the enzyme catalyst with glycolonitrile for glycolic acid production. Methods to synthesize glycolonitrile by reacting aqueous solutions of formaldehyde and hydrogen cyanide have previously been reported (U.S. Pat. No. 2,175,805; U.S. Pat. No. 2,890,238; and U.S. Pat. No. 5,187,301; Equation 3).\n\nHowever, these methods typically result in an aqueous glycolonitrile reaction product that requires significant purification (e.g., distillative purification) as many of the impurities and/or byproducts of the reaction (including excess reactive formaldehyde) may interfere with the enzymatic conversion of glycolonitrile to glycolic acid, including suppression of catalyst activity (i.e., decreased specific activity). In particular, it is well known that formaldehyde can create undesirable modifications in proteins by reacting with amino groups from N-terminal amino acid residues and the side chains of arginine, cysteine, histidine, and lysine residues (Metz et al., J. Biol. Chem., 279 (8): 6235-6243 (2004)). Suppression of catalyst activity decreases the overall productivity of the catalyst (i.e., total grams of glycolic acid formed per gram of catalyst), adding a significant cost to the overall process that may make enzymatic production economically non-viable when compared to chemical synthesis. As such, reaction conditions are needed that can help to protect the enzymatic activity against undesirable impurities that decrease the activity of the catalyst.\nA method of producing high purity glycolonitrile has been reported by subjecting the formaldehyde to a heat treatment prior to the glycolonitrile synthesis reaction (U.S. application Ser. No. 11/3143865 and U.S. application Ser. No. 11/314905; Equation 3). However, glycolonitrile can reversibly disassociate into formaldehyde and hydrogen cyanide. As such, there remains a need to protect nitrilase activity against the undesirable effects of both formaldehyde and hydrogen cyanide produced by dissociation of glycolonitrile.\nU.S. Pat. No. 5,508,181 also describes similar difficulties related to rapid enzyme catalyst inactivation when converting nitrile compounds to α-hydroxy acids. Specifically, U.S. Pat. No. 5,508,181 provides that α-hydroxy nitrile compounds partially disassociate into the corresponding aldehydes, according to the disassociation equilibrium. These aldehydes were reported to inactivate the enzyme within a short period of time by binding to the protein, thus making it difficult to obtain α-hydroxy acid or α-hydroxy amide in a high concentration with high productivity from α-hydroxy nitriles (col. 2, lines 16-29). As a solution to prevent enzyme inactivation due to accumulation of aldehydes, phosphate or hypophosphite ions were added to the reaction mixture. Similarly, U.S. Pat. No. 5,326,702 describes the use of sulfite, disulfite, or dithionite ions to sequester aldehyde and prevent enzyme inactivation, but concludes that the concentration of α-hydroxy acid produced and accumulated even by using such additives is not sufficient for most commercial purposes.\nMoreover, U.S. Pat. No. 6,037,155 teaches that low accumulation of α-hydroxy acid product is related to enzyme inactivation within a short time due to the disassociated-aldehyde accumulation. These inventors suggest that enzymatic activity is inhibited in the presence of hydrogen cyanide (Asano et al., Agricultural Biological Chemistry, Vol. 46, pages 1165-1174 (1982)) generated in the partial disassociation of the α-hydroxy nitrile in water together with the corresponding aldehyde or ketone (Mowry, David T., Chemical Reviews, Vol. 42, pages 189-283 (1948)). The inventors address the problem of aldehyde-induced enzyme inactivation by using microorganisms whose enzyme activity could be improved by adding a cyanide substance to the reaction mixture. The addition of a cyanide substance limited the disassociation of α-hydroxy nitrile to aldehyde and hydrogen cyanide. While this tactic provides a benefit to the system, it only addresses one aspect associated with enzyme inactivation in conversion of glycolonitrile to glycolic acid, in that, as stated above, glycolonitrile is known to reversibly disassociate to hydrogen cyanide and formaldehyde, and both are known to negatively effect enzyme catalyst activity.\nA separate process has been developed to protect the specific activity of an enzyme catalyst having nitrilase activity when converting glycolonitrile to glycolic acid in the presence of formaldehyde (see copending U.S. application Ser. No. 11/931,069 incorporated herein by reference), where significant improvements in catalyst activity and stability were achieved by adding an amine protectant to the reaction mixture, or by immobilization of the nitrilase catalyst in or on a matrix that is comprised of an amine protectant, e.g. PEI, polyallylamine, PVOH/polyvinylamine, etc. In that system, the specific activity of the catalyst in the presence of formaldehyde is improved, but does not address, altogether, issues related to the loss in specific activity of recovered catalyst activity in consecutive batch reactions with catalyst recycle.\nU.S. Pat. No. 4,288,552 discloses (column 1, lines 46-49, and column 2, lines 50-55) that glutaraldehyde-sensitive enzymes (such as thiol-enzymes (e.g, nitrilase) and others with an SH group in or very near the active site of the enzyme molecule) are inactivated by thiol-reactive agents such as glutaraldehyde. Therefore, use of glutaraldehyde to improve the retention of initial catalyst activity during hydrolysis of glycolonitrile to glycolic acid was heretofore unpredictable. Said unpredictable benefit is demonstrated herein.\nTherefore there is a need for a process that provides improved retention of initial nitrilase activity during hydrolysis of glycolonitrile to glycolic acid, thereby improving overall catalyst activity, catalyst productivity, and volumetric productivity for the conversion of glycolonitrile to glycolic acid."}
{"text": "This invention relates to a method and apparatus for performing capillary electrophoresis. More particularly, this invention relates to an improved method and apparatus for capillary electrophoresis that permits concentration of solutes from a solution prior to separating the solutes by capillary electrophoresis.\nCapillary electrophoresis (CE) is an efficient analytical separation technique for analysis of minute amounts of sample. CE separations are performed in a narrow diameter capillary tube, which is filled with an electrically conductive medium termed the \"carrier electrolyte.\" An electric field is applied between the two ends of the capillary tube, and species in the sample move from one electrode toward the other electrode at a rate which is dependent on the electrophoretic mobility of each species as well as on the rate of fluid movement in the tube. CE may be performed using gels or liquids, such as buffers, in the capillary. In one liquid mode, known as free zone electrophoresis, separations are based on differences in the free solution mobility of sample species. In another liquid made, micelles are used to effect separations based on differences in hydrophobicity. This is known as Micellar Electrokinetic Capillary Chromatography (MECC).\nCE is advantageous for several reasons. These include fast separation speed, high resolution and small sample size. For example, separation speeds using CE can be 10 to 20 times faster than conventional gel electrophoresis, and no post-run staining is necessary. In part, high resolution can be obtained through the use of high voltages because of the rapid dissipation of heat by the capillary. Further, band broadening is minimized due to the narrow capillary inner diameter In free-zone electrophoresis, the phenomenon of electroosmosis, or electroosmotic flow (EOF) occurs. This is a bulk flow of liquid which affects all of the sample molecules regardless of charge. Under certain conditions EOF can contribute to improved resolution and separation speed in free-zone CE.\nIn order to achieve the high resolution that CE is capable of, it is necessary that the sample be confined to a narrow starting zone when the electrophoretic process begins. This limits the volume of sample that can be introduced into the capillary to a very small fraction of the total capillary volume. Further, the lowest concentration of material that can be detected in CE by ultraviolet and visible absorbance detection is severely limited by the small inner diameter of the capillary, since absorbance detection is generally carried out with a beam of radiation that is transverse to the axis of the capillary. These facts, taken together, mean that CE has not been able to analyze samples at concentrations as low as those that can be analyzed by other techniques such as liquid chromatography. One approach to overcome this limitation has been described by several investigators. This approach involves dissolving the sample in a buffer or electrolyte having an ionic strength substantially lower than the carrier electrolyte used to carry out the electrophoretic separation. The sample so dissolved is introduced into the capillary in the normal fashion although a larger sample volume is now permitted. When the electric field is applied to begin the electrophoretic separation, the field strength in the sample zone will be higher than in the surrounding carrier electrolyte because the ionic strength and hence the conductivity is lower in the sample zone. As a result of the higher field strength in the sample zone, the ions in the sample will migrate at a higher speed than ions in the carrier electrolyte. Thus the sample ions will tend to pile up at the interface of the low conductivity sample zone and the carrier electrolyte because they will slow down once they enter the carrier electrolyte. This process effectively produces a narrow sample zone and has allowed an improvement of perhaps a factor of five in the lowest concentration of sample that can be analyzed.\nA technique for preconcentrating a sample prior to electrophoretic separation was reported by Aebersold and Morrison at the High Performance Capillary Electrophoresis Symposium in January 1990. The inner surface of a short piece of a capillary tubing was coated with an adsorptive coating. This tube was then butted up against the end of an uncoated capillary tube. This arrangement permitted a large volume of sample (which would normally have caused severely broadened peaks) to be injected into the coated length of tube and then into the uncoated tube. The sample molecules adsorbed on the coating while the solvent passed through the tube and was washed out of the capillary tube with electrolyte. A small volume of organic solvent was then introduced into the butted capillary tubes. This solvent caused the sample molecules to desorb from the coating so that they entered the uncoated tube in concentrated form. Because the organic solvent had lower conductivity than the carrier electrolyte, the previously mentioned zone narrowing then occurred. With both ends of the capillary tube in electrolyte, capillary electrophoretic separation was then performed in normal fashion. It was reported that this preconcentration step allows the use of CE with samples 5 to 10 times lower in concentration than was previously possible with CE.\nPacked capillary columns have been used in liquid chromatography to effect sample separations. Such columns have been constructed in various ways. One way of making such a column is disclosed by Jorgenson, et al in Journal of Chromatography, Vol. 218 (1981) pgs. 209-216. The column is formed by filling one end of a capillary tube with a particulate packing. The filled end then is heated in a flame to sinter the packing to form a porous frit. A slurry of chromatographic separation particles is introduced into the end of the capillary opposite to the fritted end until the capillary is filled with these particles. A second frit then is formed in the open end of the capillary in the same manner as was the first frit.\nAs disclosed in U.S. Pat. No. 4,793,920 to Cortes et al, a packed silica capillary chromatography column is produced by casting a plug of ceramic material in one end of the capillary to form a porous plug which is fused such as with heat. The chromatographic packing then is introduced into the remainder of the tube. U.S. Pat. No. 4,483,773 to Yang also discloses a means for forming a Plug and chromatographic packing within a capillary chromatographic column.\nIt would be desirable to provide a means whereby dilute sample solution can be analyzed by CE at concentrations far lower than can be analyzed by presently available CE processes."}
{"text": "This invention relates to power tools and particularly to an application mechanism for a power tool of the type comprising a housing, a motor mounted in said housing, a gearbox having an input end mounted on said motor and an output end adjacent an aperture in the housing, the gearbox being epicyclical and including a torque control ring which, when held stationery with respect to said housing, permits the gearbox to transmit torque and which, when permitted to rotate in the housing, disables torque transmission by the gearbox.\nIn such power tools it is known to dispose an application mechanism on the output end of the gearbox, which mechanism comprises a resiliently biased, adjustable actuating means mounted on said body to selectively engage said torque control clutch ring.\nAdjustment of the actuating means varies the engagement with said torque control clutch ring thereby varying the torque at which it begins to slip and at which transmission by the gearbox to its output shaft is stopped.\nIt is also known to lock the torque control ring so that it is not employed and so that no torque limitation is provided. In this instance, it is also known to dispose an application mechanism on the output end of the gearbox, which mechanism comprises a hammer arrangement whereby oscillating axial vibration can selectively be imposed on the rotary drive supplied by the gearbox.\nFinally, it is also known from DE 4038502 to provide an arrangement as first described above, i.e. a tool having a torque control mechanism mounted on the output end of the gearbox, but where a hammer mechanism is added on the front end of the torque control mechanism. Such an arrangement is versatile because it can be employed in various different ways. However, it suffers from being somewhat long having first a motor, then a gearbox, then a torque control mechanism, then a hammer mechanism, all one after the other.\nIt is also known to employ electronic torque control whereby the torque applied is fed back to an electronic control module and, if the applied torque exceeds the torque preset in the control module, power is disconnected from the motor. Thus if a hammer mechanism is on the front of the gearbox there is no penalty in terms of compactness in providing both torque control and a hammer facility but, of course, the electronics do add cost and complexity.\nIt is therefore, an object of the present invention to provide a power tool which has this versatility and yet remains compact but which is still simple and relatively inexpensive.\nIn accordance with the present invention there is provided a power tool comprising a housing, a motor mounted in said housing, a gearbox having an input end mounted on said motor and an output end adjacent an aperture in the housing, the gearbox being epicyclical and including a torque control ring which, when held stationery with respect to said housing, permits the gearbox to transmit torque and which, when permitted to rotate in the housing, disables torque transmission by the gearbox; and an application mechanism driven by said gearbox, which mechanism comprises:\na) a substantially cylindrical hollow body;\nb) a chuck drive spindle journalled in said body and defining therebetween an annular space;\nc) means limiting axial movement of said spindle in said body towards said gearbox;\nd) a rotary ratchet plate fixed on said spindle;\ne) a facing fixed ratchet plate in said body;\nf) means to engage said ratchet plates together such that on rotation of the spindle a reciprocating action is imposed thereon, said means including a lever projecting through a slot in said body;\ng) resiliently biased, adjustable actuating means mounted on said body to selectively engage said torque control clutch ring, said lever projecting through said actuating means."}
{"text": "The present invention relates to a process for the preparation of substituted dimethyl-(3-aryl-butyl)-amine compounds by means of homogeneous catalytic hydrogenation of dimethyl-(3-aryl-but-3-enyl)-amines.\nDimethyl-(3-aryl-butyl)-amine compounds have proved to be pharmaceutically active compounds exhibiting excellent analgesic activity and very good tolerability, see EP-A1-0 693 475. WO 2005/000788 A1 describes a process for the preparation of such compounds in which, in a first stage, substituted dimethyl-(3-aryl-but-3-enyl)-amine compounds are prepared by elimination of the tertiary hydroxyl group in substituted 4-dimethylamino-2-aryl-butan-2-ol compounds. These dimethyl-(3-aryl-but-3-enyl)-amine compounds are then hydrogenated in a second stage in the presence of metal catalysts. The heterogeneous hydrogenation proceeds in good yields with adequate activities. As expected, the stereoselectivity is not very pronounced. According to the examples in WO 2005/000788 A1, if two adjacent asymmetric C atoms are present diastereomeric ratios of from 2:1 to a maximum of 3:1 can be obtained for the trans diastereomer: cis diastereomer, that is to say always in favor of the trans diastereomer. The ratio is established automatically, essentially depends on the substrate, and can be influenced by the choice of reaction conditions to only a small extent."}
{"text": "The area of miniaturized microfluidic technology, also known as “lab-on-a-chip” or “micro total analysis systems”, is rapidly expanding with the promise of revolutionizing chemical analysis and offering enabling tools and technologies for the life sciences. The majority of molecules and compounds of interest in the life sciences are in the liquid phase and typically analytical measurements are used to conduct quantitative and qualitative trace analysis of these analytes in solution. However, due to the nature and origin of the materials under analysis, sample amounts and volumes are historically in short supply and these amounts are constantly decreasing. Moreover, as seen in drug discovery and drug development, including pharmacokinetic and proteomics applications, the need for better analytical measurements that require smaller amounts and volumes of sample is growing. Inherently, microfluidics are a “good fit” for the move to both smaller sample and volume requirements. In fact, a primary reason for miniaturization has been to enhance analytical performance of the device rather than reduce its physical size. Additionally, the ability to miniaturize the analysis with the use of microfluidics allows for integration of multiple separation techniques that enable parallel processing and also for the incorporation of several types of analytical measurements in a single device (sample handling, injections, 2D separations, reaction chambers etc.). Inherently, there are other benefits that accompany miniaturization, such as reduction in reagent and waste disposal, as well as, the reduction of the device footprint.\nThe first analytical miniaturized device was a gas chromatograph fabricated in silicon over 25 years ago. This device was designed for separating components in the gas phase. A decade later, components of a liquid chromatography column were fabricated on a silicon chip. Most early work related to miniaturization on silicon and the research focused on the fabrication of components such as microvalves, micropumps, and chemical sensors. None of the early systems implemented integrated electronics or electric fields for operation, but rather the silicon was used as a substrate for micromachining desired shapes/geometry.\nMost methods used in microfabrication were developed in the 1970's and 1980's in the silicon microprocessor industry. Typically, initial research developments were fabricated in silicon because of the extensive knowledge and tools available for silicon processing. This approach works satisfactorily for devices used for “dry” analyses, however, many microfluidic devices require the integration of on-chip electronics and/or the ability to apply electric fields to device. Because most applications in the life sciences involve samples contained in the liquid phase, the majority of micro analysis systems being developed/designed are for analysis of analytes contained in solution. The need for the ability to apply electric fields to the device becomes a serious issue when processing samples contained in the liquid phase on silicon substrates.\nA difficult scenario is encountered due to the opposing objectives of the micro total analysis system and the microprocessor technology used to make them. Typically, the microprocessor industry strives to keep microdevices “dry” and “clean” as liquids, moisture, and contaminates interfere with the device performance and operation. This highly contrasts the needs of micro total analysis systems where liquids and foreign substances (analytes, including salts) are deliberately introduced to the device. Again, this does not pose a problem for a silicon device that does not involve electronic or electrical field generation where only specific geometries are micromachined. However, major issues arise when electronics are incorporated in the devices and especially when potentials are applied for the generation of electric fields (semiconductor must be insulated for controlled electric field generation).\nThere is currently a move to perform chemical separations on-chip with the use of electric fields, for example applications such as CE, CEC, charged analyte manipulation, and charged solution manipulation. Because of the relatively strong fields needed for the separation process, research has moved to considering non-conductive substrates such as glass, quartz, and non-conducting polymers as opposed to the use of silicon as in the manufacture of semiconductors. This shift in materials is warranted because in order to form the electric fields, the substrates must be insulated in the desired areas.\nAlthough there are conventional techniques for insulating silicon substrates, the dielectric coatings currently available are designed for the electronics industry and operate under “dry” conditions. Much effort in the microprocessor industry has been expended on keeping devices dry or isolated from liquids. Additionally, the microprocessor industry has gone to great lengths to avoid contact of the electronic device with mobile ions such as salts due to the destructive nature they pose to dielectric coatings used to insulate the silicon used in microprocessors.\nMicrofluidic devices requiring the application of relatively high voltages and electric fields for the manipulation of liquids and samples are mainly fabricated on insulating substrates because of their insulating properties. Application of high voltages to liquids on insulators on conductor substrates often leads to shorting or drastically reduced performance and lifetime of desired electrical properties.\nAccordingly, the art needs dielectric coatings that do not degrade, but rather maintain their electrical properties when exposed to direct voltage application and high electric field strengths while in the presence of high humidity and/or direct liquid contact (wet). The art needs to overcome current coating technology limitations and provide appropriate solutions for microfluidic device applications. The art needs microfluidic devices that take advantage of the highly developed silicon processing techniques for silicon and other substrates including micromachining as well as electronic circuit integration and electric field definition. The art lacks the ability to incorporate microfluidics and electronics in the same substrate allowing for fully integrated systems."}
{"text": "1. Field of the Invention\nThe present invention pertains to a method and apparatus for selling merchandise, and particularly to a method and apparatus for selling duty free merchandise in combination with travel tickets.\n2. Description of the Related Art\nDuty free shopping is a multi-national business with substantial revenues. In the year 2000 alone, worldwide duty free sales are estimated at $21 billion. While the market for duty free goods is substantial and growing, duty free shopping is a relatively recent development, having its origins in so-called “export stores” first established in Ireland in the 1940s. In the ensuing years, most governments enacted laws authorizing and regulating duty free shopping within their borders. Initially, duty free shops offered only items that were subject to high import taxes, such as liquor and cigarettes. However, duty free merchants soon recognized the potential for selling a wide range of luxury goods to international travelers, and over time many additional items have become available as duty free merchandise, including jewelry, watches, perfume, cosmetics, cameras, electronic goods, etc. With help form world-renowned fashion houses such as Dior® and Chanel®, today duty free shops are international showcases for first class luxury goods. International travelers, in turn, have come to appreciate the quality of the shops and often take advantage of them during airport delays and even preplanned duty free shopping trips.\nInitially, most duty free operators were local merchants who secured space at airports through personal contacts. However, airports came to realize the revenue potential of such shops and began enlarging the shopping areas for international travelers. To maximize revenue, airport management eventually began issuing requests for proposals (RFP) from persons seeking to operate duty free airport shops. The response was substantial and since then the size and number of duty free shops has expanded significantly at airports throughout the world. The potential value of a particular duty free shop will, of course, depend on a number of factors, such as the total number of international passengers enplaning at that airport, the number of international routes, the frequency of flights, the airlines that utilize the facility, e.g., upscale, leisure, charter, the location of the store relative to the gate area where passengers enplane, the local prices for comparable merchandise, the prices for the goods at the usual travel destinations, the duty or tax, if any, which the passenger may have to pay at the travel destination, the projected increase or decrease in international travelers at the airport, and the length of the lease. The revenue paid by the operator of the duty free store to the airport varies widely depending on these and other factors. Oftentimes the arrangement includes a minimum rent plus a percentage of sales.\nAirlines have also sought to profit from the duty free shopping boom by selling duty free merchandise in-flight. Typically the duty free merchandise is retained on one or more trolleys which are wheeled up and down the aisles by flight personnel at preset times. The merchandise and pricing may be in a brochure shown by flight personnel to passengers expressing an interest, or it may be in a catalog at each passenger seat, or both. Travelers purchasing merchandise typically pay in cash, checks or by credit card, though checks and credit card sales carry a risk as in-flight verification that an account is in good standing is difficult.\nThe particular assortment of goods on the trolleys varies depending on the carrier and the route. For example, presently one airline's most expensive item is a Hermes scarf which sells for $265.00, whereas another airline lists a Bvlgari watch which sells for about $1,000.00 at current exchange rates. Some airlines operate their own in-flight trolleys and keep the revenue, though typically a percentage of sales is paid to flight personnel as a commission. Other airlines, particularly those in North and South America, outsource their in-flight duty free sales programs, just as they do with other services, such as meal preparation. When in-flight duty free sales are operated by a third party provider, the provider typically pays the airline a negotiated percentage of sales depending on the commission paid to the flight crew and other factors.\nSome in-flight programs offer passengers the option of pre-ordering merchandise from a special catalog of duty free merchandise. In one variation, catalogs are forwarded to international travelers sufficiently in advance of departure such that orders can be placed and fulfilled by the departure date, with the merchandise being delivered to the passenger during flight. In another variation, catalogs are presented to international travelers on their outbound flights, and ordered merchandise is awaiting them on their return flights. These programs have the advantage that the program provider can offer the traveler a wider range of merchandise than is available on the trolley, and it minimizes the amount of inventory that the provider must maintain While all these modes of duty free shopping have had various degrees of success, they also have inherent shortcomings. For example, duty free shops face the problem of late arriving passengers, international travelers spending more time clearing security, and increased pressure on airlines to meet scheduled departure times, all of which contribute to international travelers spending less time in airport duty free shops. In-flight trolleys have a limited amount of space for showcasing goods and can carry only a limited amount of merchandise. Such trolleys also face strong competition from duty free shops, which can display and inventory a much wider range of goods. Duty free catalogs forwarded to passengers before flight face passengers who may not be at the mailing address when the catalog arrives, difficulty distinguishing their catalogs from the multitude of other catalogs which flood the mails, passengers who purchase their tickets too close to flight time to receive the catalog before departure, and passengers who subsequently cancel their flights, all of which reduce sales and/or increase expenses. Catalogs delivered in-flight for product delivery upon return solve some of these problems, but this approach is useless for passengers traveling one way or those purchasing gifts for persons at the flight destination, and also fail to satisfy the traveler's desire for immediate gratification. Furthermore, all duty free shopping suffers from the perception among many travelers that duty free prices are high, which, excepting for catalogs forwarded to the traveler prior to flight, results, in part, from the traveler's inability to compare the duty free price with prevailing retail prices.\nWhile much duty free shopping targets international airline travelers, duty free shopping is also available to international travelers traveling by boat and even land. Cruise lines, for example, market duty free goods using many of the same techniques employed by the airlines. Techniques for marketing to persons traveling internationally by land (e.g., train, bus, etc.) are necessarily more limited, and typically focus on shops at the exit point from the country where the goods are sold.\nMost countries have enacted detailed laws and regulations which enable and regulate duty free shopping within their borders. Essentially, duty free shopping works as follows. A merchant desiring to sell foreign made or domestic duty free goods purchases such goods from the manufacturer, either directly or through a distributor. As long as such goods are retained for export only, they are exempt from customs duties and as well as federal, state and local taxes. The shops from which merchants sell such goods are known as duty free or tax free shops, however both are the same, the former referring to the exemption from customs' duties upon export, and the latter referring to the exemption from government taxes imposed by the country in which the shop is situated. As used herein, “duty free” refers to sales which are exempt from customs' duties and/or taxes.\nFor example, the merchant may purchase foreign-made ties, jewelry, watches, liquor, handbags, etc. and then import those goods into the country from which the duty free sales will be made. Upon export following sale, no duty or tax is paid by the merchant on such imported goods provided, prior to sale, the goods are maintained in a bonded facility sanctioned for that purpose by the government and provided the merchant complies with all applicable regulations governing such facilities. Alternatively, or in addition, the merchant may purchase domestic goods. While domestic goods would not, in any event, be subject to duty upon sale, they are typically taxed when sold. However, provided they are maintained in a bonded facility of the type described above, domestic goods are exempt from tax upon export, which is particularly attractive in the case of goods which command a luxury tax when sold domestically, such as alcohol and cigarettes. In fact, merchants selling domestic goods duty free can typically acquire such goods from manufacturers at prices below, and sometimes well below, the prices at which the manufacturers sell such goods to conventional retail outlets, which makes such sales particularly attractive for duty free merchants. Typically, government regulations will govern the manner is which goods may be brought into and taken out of bonded facilities, and will include strict inventory and record keeping requirements to insure that duty free goods in the bonded facility do not find their way into local commerce.\nTo maintain duty free status, duty free merchandise may only be delivered to travelers departing the country from which the goods are sold for use and/or consumption outside that country, though such goods may be reimported without duty or tax, provided their value does not exceed the entry exemption allowed for imported goods. To insure that duty free merchandise is exported, most countries have regulations governing where duty free shops may be located and where duty free merchandise may be delivered to the traveler. Although a duty or tax may be due to the destination country upon arrival, many countries exempt a predetermined monetary value of imported goods intended for personal use, i.e., not intended for resale. Consequently, if the traveler purchases duty free goods for personal use under the limit set by the destination country, no duty or tax will be payable either to the country where the goods were purchased or to the destination country. And even if the purchases exceed the exemption allowed by the destination country, a duty and/or tax is only due on the overage.\nIn the United States, jurisdiction over duty free sales is relegated to the Treasury Department and implemented by the United States Customs Service, which, in turn, has authority to relegate certain responsibilities to port directors at airports and seaports. The current regulations governing duty free facilities and delivery of duty free merchandise to international travelers is found in Title 19 of the Code of Federal Regulations, §§19.35 et seq. These regulations provide, among other things, that “conditionally duty-free merchandise” (defined as merchandise sold by a duty free store on which duties and/or internal revenue taxes, if any are applicable, have not been paid) for export at airport, seaport or land exit points may only be sold and delivered to persons displaying tickets or other proof of imminent departure from the country. To insure that duty free merchandise in fact leaves the country, the regulations also include strict limitations on where duty free merchandise may be delivered to international travelers, with such delivery generally limited to at or beyond the “exit point”, which is defined in the regulations as an area in close proximity to an actual exit for departing the country. In the case of an airport the regulations provide that the exit point may be a gate holding area if the gate holding area is sufficiently secure that there is a reasonable assurance that conditionally duty-free merchandise delivered there will be exported from the country. In the case of a seaport or land border, the exit point is defined in the regulations as the point at which a departing individual has no practical alternative but to continue on to a foreign country or to return to the country of departure by returning through a U.S. Customs inspection facility. The regulations also provide that conditionally duty-free merchandise brought back into the country is subject to U.S. duty and tax, subject, however, to the personal exemption afforded international travelers arriving in the United States, which is currently $400.00 per person, though this exemption can be pooled among multiple travelers for more expensive goods and in some cases, such as goods imported from the U.S. Virgin Islands, the exemption is higher. Most other countries have restrictions on duty free sales similar to those enacted in the United States.\nIt will be apparent, therefore, that new and creative methods of generating duty free sales are limited by the restrictive government regulations governing such sales, and also by the above-discussed limitations inherent in the current methods for effecting such sales."}
{"text": "1. Field of the Invention\nThe present invention relates to an endoscope having an insertion portion provided with a first and a second bendable sections and a method for inserting the endoscope into a colon.\n2. Description of the Related Art\nAn endoscope in the medical field can be used by inserting an elongated insertion portion into a body cavity for observation of organs and the like in the body cavity or performing various treatments using a treatment instrument inserted into a treatment instrument insertion channel as necessary. In this endoscope insertion portion, a tip end (constituting) portion, a bendable section and a flexible tube portion in order from the tip end are disposed.\nWhen the insertion portion of the endoscope is to be inserted into a body cavity, an operator grips the flexible tube portion and operates an operation knob disposed at an operation portion of the endoscope by a predetermined amount so as to bend the bendable section in a desired direction while pushing it into the body cavity. This type of insertion portion of the endoscope has various improvements to reinforce insertion ability into a curved body cavity.\nFor example, in the insertion portion of the endoscope described in Japanese Utility Model Publication No. 1-22641, a first bendable section which can be actively operated from outside to be bent in four directions and a second bendable section which is provided with a stay coil and a joint ring and is very easily bent and in a structure passively bent in four directions are provided consecutively in order from the tip end side.\nAlso, Japanese Unexamined Patent Application Publication No. 2002-143084 discloses an endoscope provided with a first bendable section and a second bendable section with hardness larger than that of the first bendable section."}
{"text": "Oil wells and natural gas are generally drilled into the earth and the underground oil or gas deposits are forced to the surface along the well bore by the underground pressure, or, pumped up using one or more pumps (often in a series). The well-bore is drilled from an oil-rig on the surface of earth using a rotating drilling bit. When the rotating drilling bit is driven into the earth, it cuts through layers of soil and rocks using a continuous flow of compressed drilling fluid (also known as “drilling mud”) supplied through a conduit, which can be coiled tubing or a drill string (composed of a contiguous series of pipes).\nDuring drilling, hollow metallic tubes (also known as “casings”) are inserted within the drilled bore to prevent the walls of bore from collapsing. In a deep enough bore, multiple hollow casings are installed vertically one above the other by screwing ends of adjacent sections with each other, thus forming a “bore casing.” For reliable installation of bore casings (and for proper functioning of equipment set-up within it), the well-bore should be kept clean by efficiently removing mud, rock debris and dirt generated during drilling, and it should also be periodically cleaned to maintain reliable functioning of the in-well equipment.\nVarious tools and methods have been proposed for cleaning of an oil well-bore. For example, U.S. Pat. Nos. 7,686,102 and 8,151,908 (both incorporated by reference) disclose well-bore cleaning devices which use jets of pressurized fluid ejecting from one end of a tool to dislodge mud, dirt or rock debris in an oil well-bore. While drilling the oil well-bore, these tools are not very efficient in removing mud/dirt stuck on inner walls of the bore casing at positions already passed by the end of the tool from which jets of pressurized fluid are ejected. So in order to clean such portions, the tools may be required to pass multiple times up and down the oil well-bore. Hence, there's need for an improved cleaning device which more efficiently and effectively cleans the base and internal walls of a well-bore without multiple passes up and down being required."}
{"text": "There are communication devices that are capable of performing proximity wireless communication. An example of such device is disclosed in JP-A-2005-045557 (counterpart U.S. publications are: US 2005/0064814 A1 and US 2007/0232232 A1). When information is transmitted to and received from another device, the communication device described in the publication JP-A-2005-045557 determines a transmission/reception direction by an operation button and transmits/receives the information. This information has a relatively small information amount and information contents displayed on the other device are directly transmitted and received.\nThe communication device described in the publication JP-A-2005-045557 has a special operation button for determining an information transmission/reception direction, which is operated by the user. In the case where a plurality of contents of a relatively large amount of information are stored, a content list needs to be first displayed before displaying the contents. However, in this case, a user interface, which sets one of the communication device and the other device to display the content list, is not disclosed."}
{"text": "The present invention relates to a titanium article having a protective coating and a method of applying a protective coating to a titanium article, particularly to a titanium aluminide article having a protective coating and a method of applying a protective coating to a titanium aluminide article.\nTitanium aluminide alloys have potential for use in gas turbine engines, particularly for turbine blades and turbine vanes in the low pressure turbine and compressor blades and vanes in the high pressure compressor and the combustion chamber diffuser section. The gamma titanium aluminides provide a weight reduction compared to the alloys currently used for these purposes.\nHowever, titanium aluminide alloys and gamma titanium aluminide alloys will require environmental protective coatings, above a certain temperature, in a similar manner to conventional nickel base alloys or cobalt base alloys.\nConvention environmental protective coatings for nickel base alloys and cobalt base alloys include aluminide coatings, platinum coatings, chromium coatings, MCrAlY coatings, silicide coatings, platinum modified aluminide coatings, chromium modified aluminide coatings, platinum and chromium modified aluminide coatings, silicide modified aluminide coatings, platinum and silicide modified aluminide coatings and platinum, silicide and chromium modified aluminide coatings etc. Aluminide coatings are generally applied by the well known pack aluminising, out of pack, vapour, aluminising or slurry aluminising processes. Platinum coatings are generally applied by electroplating or sputtering. Chromium coatings are generally applied by pack chromising or vapour chromising. Silicide coatings are generally applied by slurry aluminising. MCrAlY coatings are generally applied by plasma spraying or electron beam physical vapour deposition.\nThermal barrier coatings include yttria stabilised zirconia and magnesia stabilised zirconia etc. Thermal barrier coatings are generally applied by plasma spraying or electron beam physical vapour deposition.\nHowever, these conventional protective coatings are not as adherent to titanium aluminide alloys in particular, or titanium alloys in general, as they are to nickel base alloys or cobalt base alloys. This is due, we believe, to the titanium oxide formed on the titanium aluminide or titanium alloy.\nAccordingly the present invention seeks to provide a novel protective coating for a titanium article and a novel method of applying a protective coating to a titanium article.\nAccordingly the present invention provides a titanium alloy article having a protective coating on the titanium alloy article, the protective coating comprising a coating of silicate glass having a chromium oxide filler.\nPreferably the protective coating comprises an oxide layer on the titanium alloy article and the coating of silicate glass having the chromium oxide filler on the oxide layer.\nPreferably the titanium alloy article comprises a titanium aluminide, more preferably the titanium alloy article comprises a gamma titanium aluminide.\nPreferably the oxide layer comprises titanium oxide.\nPreferably the protective coating comprises a boron titanate silicate glass having a chromium oxide filler.\nPreferably the titanium alloy article comprises a turbine blade, a turbine vane, a compressor blade, or a compressor vane.\nThe present invention also provides a method of applying a protective coating to a titanium alloy article comprising depositing a coating comprising a silicate glass having a chromium oxide filler.\nPreferably the method comprises forming an oxide layer on the titanium alloy article and depositing the coating comprising silicate glass having a chromium oxide filler on the oxide layer.\nPreferably the titanium alloy article comprises a titanium aluminide, more preferably the titanium alloy article comprises a gamma titanium aluminide.\nPreferably the oxide layer comprises titanium oxide.\nPreferably the protective coating comprises a boron titanate silicate glass having a chromium oxide filler.\nPreferably the titanium alloy article comprises a turbine blade, a turbine vane, a compressor blade, or a compressor vane.\nPreferably the method comprises depositing the boron titanate glass and chromium oxide filler by spraying with a binder.\nPreferably the method comprises drying the protective coating, heating the protective coating at 100xc2x0 C. for 1 hour and heating the protective coating at 1030xc2x0 C. for 10 to 20 minutes to fuse the protective coating."}
{"text": "Allogeneic hematopoietic cell transplantation (HCT) is a curative therapy for hematological malignancies as well as refractory autoimmune diseases such as systemic lupus erythematosus (SLE) (Shizuru 2004; Sullivan 2004; Sykes 2005). HCT results in allogeneic hematopoietic chimerism. For subjects with hematological malignancies, HCT allows for administration of higher and more potentially effective dosages of chemotherapy and radiotherapy. For subjects with abnormal nonmalignant lymphohematopoietic systems, HCT allows for replacement of the abnormal lymphohematopoietic system with a healthy one. HCT is often referred to as bone marrow transplantation (BMT). However, hematopoietic stem cells for use in HCT can be collected from bone marrow, peripheral blood, or umbilical cord blood.\nIn classical HCT, recipients are conditioned with total body irradiation (TBI) and/or high dose chemotherapy. However, both TBI and chemotherapy are highly toxic, even at reduced intensity. In addition, TBI and chemotherapy play a critical role in initiating graft versus host disease (GVHD) (Ferrara 2004). Toxicity and GVHD have limited the application of HCT. Therefore, there is a need for improved methods of conditioning a recipient for HCT."}
{"text": "The present invention relates to aerial lift platform apparatus, and more particularly to such apparatus which is mounted on a chassis having endless tracks or treads which engage the ground.\nAerial lift platform apparatus have been provided in the past, and their use has recently been significantly increasing. Such apparatus is in the form of a self-propelled vehicle having a chassis which is mounted on ground engaging wheels. Smaller units often have three wheels, while larger units, having a vertical reach of between approximately forty feet and one hundred and twenty-five feet, are mounted on four wheels. The chassis supports a rotatable upper works, which is capable of rotating through 360 degrees. A boom is carried by the upper works, having at or adjacent its outer end a workman's platform, usually provided with upstanding walls so as to form a basket-like structure, for purposes of safety. The boom is often extensible, and is caused to rotate through 360 degrees, due to its mounting on the upper works. In addition, the boom may be raised and lowered about a horizontal pivot. The upper works usually contains a prime mover, such as an internal combustion engine, which drives one or more pumps, to thereby deliver power to rotary motors positioned at one or more of the wheels in order to propel the vehicle, and also to deliver hydraulic fluid under pressure to various cylinders, such as for steering and boom extension and retraction. Rotation of the upper works is also effected by a hydraulic motor driven by the pump. All controls are provided in duplicate, so that there is a set of controls accessible to a person standing by the machine, and also a set of controls is provided at the platform, so as to be accessible to the operator or workman in the platform.\nAn aerial lift platform has been provided, mounted on a tracked vehicle chassis, instead of wheels. There was thereby gained the advantage of permitting operation where the supporting surface was relatively soft, and could not sustain a vehicle of such weight when mounted on wheels. This construction was a known crawler type construction which provided a steel driving sprocket and a steel idler at the ends, and intermediate steel rollers which were rigidly mounted to the chassis. This construction, while suitable for use in terrain of a relatively soft nature, provided no shock absorbing or cushioning for the operator, who could only drive the vehicle from the the operator's platform.\nIn another embodiment, an aerial lift platform was provided which was mounted on tracks, and was provided with a cab on the vehicle, which cab had controls for moving the vehicle. The operator's platform was not provided with controls for moving the vehicle, so that the operator, when the vehicle was moving, could only be in the cab which was mounted directly on the vehicle. Once the aerial lift platform was set in position, an operator occupied the platform, and by manipulation of the controls in the platform on the boom, was able to rotate the upper works, raise and lower the boom and extend and retract the boom, but could not, as noted, cause the vehicle to move. In this construction, pneumatic tires were provided, for both drive and intermediate or idler wheels, and the idler wheels were mounted on independent spring loaded axles. That is, one or more of the idlers on one side of the chassis could be moved upwardly, or downwardly, relative to a reference point on the chassis, without causing movement of any idler on the opposite side, relative to the same or a corresponding reference point on the chassis.\nThe aerial lift platform apparatus having rigid suspension also had a cab on the chassis, and thus required two operator's stations, one at the cab and one at the platform, and there was thereby required extra cost; in addition there was the problem that the vehicle could not be moved by the operator in the cab and thus was more expensive from the point of view of personnel cost, and was somewhat less efficient. On the other hand, where the shock absorbing or cushioning type suspension was utilized, including pneumatic tires and spring loaded intermediate wheels, stability was not as great as desired, particularly when the boom was extended over one side. It has been found that with the boom extended over one side, there is produced an over-turning moment tending to tilt the machine and overturn it to the side over which the boom extends. The problem is known to be more severe where there is a shock absorbing or cushioning or so-called \"sprung\" chassis, and this has been solved in the past, for instance in connection with truck mounted cranes, by the use of outriggers. Thus, the truck mounted cranes have the outriggers in stored, non-use position during movement of the truck mounted cranes from one place to another, but when relatively heavy loads are to be lifted, the outriggers are utilized, and their positioning in effect renders the spring mounting of the chassis on the ground engaging wheels inoperative. That is, during the lifting of loads near the capacity of the truck mounted crane, the spring suspension of the chassis is made inoperative and the supporting of the chassis is by the outriggers, which bypass the spring suspension of the chassis. Of course, with such cranes, it is not possible to move them when the outriggers are positioned.\nThere are known, also, mobile cranes which are intended to lift and move a load by movement of the entire truck mounted crane. In such units, however, although outriggers are not used, movement of the entire truck mounted crane is effected with the load and the boom centered, and not extending over the side.\nVarious vehicles are known which utilize tracks, and have pneumatic tires, together with spring-loaded independently movable intermediate or idler wheels. In these constructions, however, there is no boom, and therefore no overturning moment resulting from a boom positioned over the side."}
{"text": "A conventional wrench with a box end is used for a long time and the box end includes a toothed inner periphery so as to engage a polygonal object such as a nut. The object is enclosed by the box end and cannot be slid from the box end during rotation of the wrench. Nevertheless, if there is an obstacle located beside the object, the handle of the wrench cannot be rotated freely, the user has to remove the box end from the nut completely and then re-engage the nut after the handle is arranged to a proper position. This increases the time required to tighten or loosen the nut and is not an efficient tool. To improve the problem, a ratchet box end tool is developed and allows the user to reciprocally operate the box end wrench without removing the box end from the nut. However, the ratchet wrench involves too many parts and the assembling processes are complicated.\nThe present invention intends to provide a box end wrench without ratchet mechanism and can be used as a ratchet wrench."}
{"text": "A common type of human injury occurs when a finger or toe is smashed, such as by an inadvertent hammer blow, door slam, etc. A result of these accidents is a very painful bruise and swelling under the nail, called a subungual hematoma.\nIn order to relieve the swelling under the nail, the nail is perforated to permit blood or other body fluids to ooze out to relieve pressure. By directing the pooling blood out of the restricted space under the nail, several beneficial results occur. First, the pressure under the nail is relieved so that the swelling and throbbing pain is reduced. Second, discoloration of the nail is lessened. Finally, premature loss of the nail is avoided. In the prior art, the nail has been perforated by burning the nail with a red-hot wire, or by puncturing with the tip of a hypodermic needle. A problem is lack of control: the delicate tissue under the nail can also be perforated, which can lead to additional pain and worsening of the wound.\nNail boring devices (called trephinators) are known in the art. U.S. Pat. No. 6,015,418 and US Patent Application 2006/0225757 describe hand drills for drilling a pressure-relief hole in the nail (2006/0225757 also has an electrically powered drill). However, there is no provision of making sure the drill does not inadvertently puncture the tissue under the nail.\nU.S. Pat. No. 5,645,554 attempts to provide depth controlling means so that the drill does not inadvertently puncture the tissue under the nail. The depth controlling means is a tube which is added to a standard drill bit. It is placed around the bit and has an internal circular side wall, which is engaged with one or more flutes of the drill. The enables the tube to move longitudinally along the drill as the tube rotates on the drill."}
{"text": "As metal oxide semiconductor field effect transistors (MOSFETs) scale down in size, higher performance can be achieved by bringing metal silicide contacts closer to the gate conductor. There are difficulties however in forming self-aligned metal silicide contacts close to the gate conductor edge. For example, relatively thick silicides (on the order of about 20 nm or greater), which are required to meet sheet resistivity in the diffusion and polycide, will consume silicon and can interfere with the integrity of an ultra shallow junction that is typically present near the gate edge under the spacer.\nFor ultra-thin body MOSFETs in which the device channel has a thickness of about 20 nm or less, the thickness of the silicon available to form a silicide contact is limited. So-called raised source/drain regions (RSD, or also called elevated source/drain regions) can mitigate this problem. However, RSD regions that are positioned close to the gate conductor edge will increase the capacitance between the gate and the source/drain regions.\nIn view of the above, there is a need for providing a MOSFET structure having self-aligned metal suicide contacts that are close to the gate conductor edge that do not consume sufficient silicon such that the integrity of the ultra shallow junction is not effected. Moreover, a MOSFET structure is needed in which the silicide located in proximity to the gate conductor edge does not increase the capacitance between the gate and the source/drain regions."}
{"text": "This invention relates to an optical disk recording thereon compressed data of an image, and to an optical disk reproduction apparatus for reproducing the image data from the optical disk.\nA so-called xe2x80x9cCD-ROMxe2x80x9d is a typical example of those systems which reproduce digital data by using an optical disk. The CD-ROM records data for a computer on an optical disk having the same physical format as CDs for audio use, and has the data format to be next described. A data string recorded on the optical disk comprises the smallest unit referred to as a xe2x80x9cframexe2x80x9d, and each frame contains digital data such as sync data, subcode, main information and an error correction code.\nFurther, the CD-ROM employs the sector structure in which 98 frames (2,352 bytes) are gathered into 1 sector, and each sector comprises a 12-byte sync data, a 4-byte header data representing an address and a mode, a 2,048-byte digital data and a 288-byte error detection/correction code.\nOn the other hand, a system comprising the combination of an inter-frame prediction with orthogonal transform, quantitization and variable-length encoding is well known as an encoding system of moving picture, and an MPEG system of the ISO (International Organization of Standardization) is also based on this system. In the case of MPEG2, for example, the bit stream of the encoded image data is divided into six hierarchical layers, i.e. a sequence layer, a GOP (Group of Pictures) layer, a picture layer, a slice layer, a macro-block layer and a block layer. Among them, the GOP layer contains three kinds of data, that is, an I picture encoded from the information alone without using inter-frame prediction, a P picture generated by executing prediction from the I picture or P picture and a B picture generated by bidirectional prediction. The sequence layer comprises a GOP containing image data starting from the I picture and obtained by gathering the I picture, the P picture and B picture into one group, and an SH (Sequence Header) added to the leading part of the GOP.\nWhen moving picture are converted to compressed image data by high efficiency encoding, there is known a system which encodes the signal by reducing a compression ratio for scenes having vigorous motion, or in other words, at a high transfer rate, and by increasing the compression ratio for scenes having small motion, or in other words, at a low transfer rate. The variable transfer rate compressed image data so encoded can reduce image deterioration due to compression in comparison with compressed image data of a fixed transfer rate obtained by fixing the compression ratio at a mean value.\nAn apparatus for recording the compressed image data having the variable transfer rate or the fixed transfer rate into the optical disk such as the CD-ROM and reproducing the data has been announced already.\nThe prior art technology described above does not particularly consider so-called xe2x80x9ctrick playxe2x80x9d such as n-time speed variable speed reproduction with respect to standard speed reproduction or reverse reproduction of the image data recorded on the optical disk. In reproduction of the image data, for example, n-time speed variable speed reproduction or reverse reproduction is generally required besides continuous reproduction of a standard speed, and a reproduction apparatus which satisfies these requirements becomes necessary. High speed retrieval is also required for the retrieval operation.\nIt is an object of the present invention to provide an optical disk and an optical disk reproduction apparatus each of which can conduct various kinds of trick play of the optical disk recording thereon compressed image data, and can make a retrieval operation at a high speed.\nTo accomplish the object described above, the optical disk according to the present invention records specific information necessary for trick play into an arbitrary area such as a TOC (Table of Contents) or a leading sector (sector 0) of the optical disk, and adds a sector address to each sector.\nThe optical disk reproduction apparatus according to the present invention includes means for extracting and reproducing an I picture, a P picture and a B picture contained in a GOP layer inside a bit stream of compressed image data by looking up a trick play table.\nAfter the optical disk is loaded to the optical disk reproduction apparatus, a system microcomputer first reads the trick play table recorded on the optical disk and stores it into a work area. When trick play is effected, the address of the sector to be read out is determined by looking up the necessary trick play table, and the address is retrieved on the optical disk so as to reproduce an image.\nBecause the address of the sector to be read out is determined by looking up the trick play table during trick play other than normal reproduction, trick play can be easily executed, and a retrieved reproduction image can be quickly obtained in the retrieval operation."}
{"text": "Worldwide demand for cooling and refrigeration systems are rising exponentially. Driven by a warming planet and a rapidly expanding middle class in developing economies, the use of cooling systems (e.g., air conditioners, cold storage units, beverage coolers, bulk milk coolers, etc.) is surging. In most countries, energy consumed by cooling systems account for a substantial portion of the total power used. Almost all current cooling systems operate on the vapor-compression thermodynamic cycle, where a circulating liquid refrigerant absorbs and removes heat from a cooled space (such as, for example, an enclosed room, cabinet, etc.), and rejects the heat elsewhere. However, such cooling systems consume large amounts of power. A recent intergovernmental study estimates that power consumption for residential air conditioning alone will increase over thirty-fold by the year 2100.\nFurthermore, cooling systems (such as air conditioners) in a geographic area tend to turn on at the same time causing a surge in power use. To accommodate this demand surge, utilities in developed countries use “pecker plants” (gas turbines, etc.) which can provide instantaneous power at a significantly higher cost. Since consumers are unwilling to pay the higher rates, many parts of the developing world encounter blackouts as the demand exceeds supply. To reduce the cost of peak power, and the inconvenience of blackouts, it is desirable to shift some of the high electricity loads to non-peak time (e.g., morning, night, etc.). One way to accomplish this is by storing the energy needed to operate the cooling systems in batteries. However, electric battery cost is very high. Embodiments of the current disclosure provide systems and methods to alleviate some of these deficiencies. The scope of the current disclosure, however, is defined by the attached claims, and not by the ability to solve any specific problem."}
{"text": "A security dongle is a piece of hardware to be used in physical connection with a computer as a security mechanism for a piece of software on the computer. The security dongle can be used for copy prevention and/or software feature enablement. Typically, the security dongle authenticates the piece of software based on a unique identifier of the security dongle. Without the security dongle, the piece of software will run in a restricted mode or not run at all."}
{"text": "Power storage technology is an important technology for efficient use over the whole fields of energy such as efficiency of power use, enhancing ability or reliability of power supply systems, expanding introduction of new and renewable energy having large fluctuation depending on time, and energy regeneration of moving bodies, and its potential and demands for social contribution have gradually increased.\nStudies on secondary batteries have been actively conducted in order for adjusting the balance of demand and supply in semi-autonomous regional power supply systems such as micro-grid, properly distributing non-uniform outputs of new and renewable energy generation such as wind power or solar power generation, and controlling affects such as voltage and frequency changes occurring from differences with existing power systems, and expectations for the use of secondary batteries have increased in such fields.\nParticularly, when examining properties required for secondary batteries used for high capacity power storage, energy storage density needs to be high, and as high capacity and high efficiency secondary batteries suited for such properties, redox flow batteries (RFB) have recently received attention.\nLike general secondary batteries, redox flow batteries store electric energy inputted through a charging process after converting to chemical energy, and output the stored chemical energy after converting to electric energy through a discharge process. However, such redox flow batteries are different from general secondary batteries in that electrode active materials storing energy are present in a liquid state instead of a solid state, and therefore, a tank or a storage container storing the electrode active materials is required.\nAs described above, redox flow batteries have properties such as capable of being manufactured to have large capacity, requiring lower maintenance costs, operable at room temperature, and capable of designing capacity and output each independently, and therefore, considerable studies on redox flow batteries as high capacity secondary batteries have been currently conducted.\nAmong these, vanadium redox flow batteries using vanadium ions have received attention as a next generation energy storage device, however, there are problems of capacity decrease in the redox flow battery due to phenomena such as separator (or ion-exchange membrane) cross-over of vanadium ions, generation of hydrogen in an anode, and oxidation reaction of vanadium ions when exposed to air, and studies for improving these have been consistently ongoing.\nAmong these, a separator cross-over phenomenon of vanadium ions is caused by ion imbalance between an anode electrolyte and a cathode electrolyte occurring by the anode electrolyte and the cathode electrolyte having different oxidation numbers, and as a result, a problem of causing deterioration in the battery capacity occurs.\nSpecifically, V2+ and V3+ ions of the anode electrolyte permeate the membrane in a relatively higher rate compared to V5+ and V4+ ions of the cathode electrolyte, and as a cycle progresses, the vanadium ion concentration rapidly increases in the cathode electrolyte. Accordingly, the concentration of vanadium ions decreases in the anode electrolyte, and as a result, ion balance between the anode and cathode electrolytes is broken causing a problem of cycle deterioration when operating a battery."}
{"text": "Typical fluorescent lamps contain electrodes that, when operated at the lamp's rated discharge current, are heated to cause thermionic emission of electrons. Such lamps typically contain two electrodes. Each electrode serves as an anode during a first half-cycle of the alternating current (AC) provided to the electrodes, while the other electrode serves as cathode. During the subsequent half-cycle, the electrodes swap roles. Thus, each electrode serves as cathode and anode on alternating half-cycles.\nSuch fluorescent lamps may be dimmed by reducing the current supplied to the electrodes. Reducing the current, however, also reduces the electrode temperature. If the cathode, for example, is not sufficiently hot, it may not have sufficient thermionic emission to maintain the discharge as a thermionic arc. Rather, it may be forced to operate in a “cold-cathode” (i.e., high cathode-fall, high sputtering) mode. The resulting sputtering damage may cause the cathode, and thus the lamp, to fail within a few hours. Less destructive, but equally fatal in the long run, is the fact that a not-quite-hot-enough cathode will have a higher-than-normal cathode fall even though the cathode continues to operate in the thermionic-arc mode. If the cathode fall exceeds the so-called “disintegration voltage” or “sputter voltage,” incoming mercury ions bombard the cathode with sufficient energy to dislodge (i.e., “sputter”) surface atoms, bringing about increased rate-of-loss of cathode coating and short life.\nIn order to avoid these effects in fluorescent lamp dimming, an auxiliary electrode heating current may be supplied to the electrode filament to heat the filament sufficiently, by Ohmic heating, to cause thermionic emission. At low dimming currents with minimal heating of the cathode from the discharge current, the auxiliary supply may be the only heat source available to maintain cathode temperature. The auxiliary supply may be a low voltage, typically <6 volts, AC supply connected to the two ends of the filament structure holding the emissive coating. Resistive heating in the filament then furnishes the necessary heating power to maintain the cathode temperature at a desired level. The heating level and corresponding cathode temperature may be controlled by adjustments in the voltage or current furnished by the auxiliary heating power supply.\nObviously, the lower the dimming level at which the lamp is operated, the higher the auxiliary electrode heating current that will be required. If the voltage of the auxiliary heat supply is too low, then the cathode is too cold, the thermionic-arc cathode fall is too high, sputtering occurs with accelerated loss of cathode coating, and short lamp life results. Accordingly, it would be desirable to define a lower limit for the auxiliary electrode heating current in order to keep the lamp life within a reasonable range. On the other hand, if the voltage of the auxiliary heat supply is too high, the cathode temperature is too high, and excessive evaporation of cathode coating leads to short lamp life, even though the cathode fall is maintained well below the disintegration voltage. Steering a course between the Scylla and Charybdis perils represented by sputtering or excess evaporation is therefore desirable in selecting an appropriate cathode-heating-voltage profile as a function of discharge current. Such considerations may be particularly useful in the design of dimming ballasts.\nBy techniques known in the art, cathode temperature may be measured with an optical pyrometer, provided special lamps with phosphor wiped away from the ends are used to render the cathode visible. Alternatively, average cathode temperature may be determined in lamps without wiped ends by measuring the ratio of hot-to-cold-resistance of the cathode tungsten coil.\nA technique is known for determining heater current as a function of discharge current in one particular lamp type of one particular wattage (see F. S. Ligthart, H. Ter Heyden, and L. Kastelein (Paper 17L, 5th International Symposium on Science and Technology of Light Sources, York, England 1989). This technique, however, involves life-testing a number of lamps at various values of heater current and discharge current for a long period of time. Examination of the resulting discolorations on the ends of the lamps could discriminate between sputtering, which may lead to so-called “end band” discoloration, excess evaporation, which may lead to so-called “cathode spotting” discoloration, and satisfactory heater voltage, which exhibits little or no discoloration.\nFIG. 1 provides typical sputter and vaporization curves obtained using this technique. Heater-current versus lamp-current points where excessive evaporation was found are identified with a □; points where excessive sputtering was found are identified with a o; points where neither excessive evaporation nor excessive sputtering were found are marked with an x. Points lying between the two curves may be considered acceptable. Points lying below the sputter curve may lead to sputtering. Points lying above the vaporization curve may lead to excessive evaporation.\nThough such a technique may provide information that is useful in determining the correct heater-current profile versus dimming current, it is far too cumbersome for an electronic ballast manufacturer to employ efficiently. To provide comprehensive results, it would have to be repeated for every different wattage of every different lamp type. In addition, cathode designs employed by different lamp manufacturers for the same lamp type are different, requiring testing of a number of different versions of the same lamp type.\nFIG. 2 provides a plot of cathode fall as a function of phase angle for a typical fluorescent lamp. Specifically, FIG. 2 shows cathode fall as a function of phase angle in a typical T12 Rapid Start fluorescent lamp operating at rated current and heater voltage. Measurements of cathode fall were made using special lamps equipped with so-called “Langmuir Probes” (see John F. Waymouth, “Electric Discharge Lamps,” MIT Press 1971, Chapter IV and Appendix B). It should be understood that, as the data provided in FIG. 2 is on an absolute, rather than relative, basis, the plot of FIG. 2 provides a standard against which other methods for measuring cathode fall may be compared.\nFIG. 3A provides a plot of cathode and anode falls for a typical fluorescent lamp. FIG. 3B provides a plot of arc current supplied to produce the cathode and anode falls plotted in FIG. 3A. Cathode fall in 60-Hz AC operated fluorescent lamps was measured via a method, attributable to Hammer, et al., wherein a capacitive probe, which may be a foil sheet, for example, is wrapped around a portion of the lamp that contains the electrode (see, for example, E. E. Hammer, “Comparative Starting Characteristics in Typical F40 Systems”, Preprint, IESNA Conference Minneapolis Minn. 1988, and its published version, JIES Winter 1989, p 64; and E. E. Hammer, “Cathode Fall Relationships in Fluorescent Lamps”, Preprint #69, IESNA Conference, Miami Beach Fla., 1994, and its published version, JIES, Winter 1995, p116). The probe picks up fluctuations of plasma potential in proximity to the electrode, and presents them to an oscilloscope for detection. As the negative glow in front of the cathode is approximately an equipotential blob of high-density plasma at a potential (positive relative to the cathode) that is equal to the cathode fall, fluctuations of plasma potential on the cathode half cycle may be interpreted as the signature of fluctuations of cathode fall during the half cycle.\nPositive swings of potential are attributed to cathode fall (negative glow plasma that is positive with respect to the electrode) while negative swings are identified with anode fall (negative glow plasma that is negative with respect to the electrode). When allowance is made for the difference in lamp-current waveform, the shape of the curve agrees well with that shown in FIG. 2. The jagged fluctuations seen in FIG. 3A on the anode half cycle are so-called “anode oscillations.”\nBecause capacitive coupling causes a loss of DC reference, the Hammer method provides no information as to the value of the zero of potential. Thus, although the Hammer method may provide qualitative information about the shape of the cathode and anode fall waveforms, and about the peak-to-peak difference between peak cathode fall and peak anode fall, it does not provide the magnitude of either. Further, fluctuating potential signals from the cathode heater voltage may be picked up. Also, a very high input impedance is required for the measuring oscilloscope, which precludes the use of the Hammer method on small-diameter, compact fluorescent lamps. Additionally, capacitance between the signal lead and the shielded cable shunts the signal to ground, which reduces the apparent amplitude of the cathode fall variation. It would be desirable, therefore, if apparatus and methods were available for making capacitive measurements of the magnitude of cathode fall in fluorescent lamps, without the limitations exhibited by the Hammer method."}
{"text": "The quantitative determination of analytes in biological fluids is useful in the diagnosis and treatment of physiological abnormalities. For example, determining the glucose level in biological fluids, such as blood, is important to diabetic individuals who must frequently check their blood glucose level to regulate their diets and/or medication.\nElectrochemical methods have been used for such purposes. An electrochemical biosensor may use an analyte specific enzyme, such as glucose oxidase or glucose dehydrogenase, to catalyze the oxidation of glucose in a whole blood sample. During the catalytic oxidation by the enzyme, the redox center of the enzyme accepts the electrons from the analyte.\nThis redox center could be the flavin adenine dinucleotide (FAD) of glucose oxidase, or the enzyme's cofactor such as pyrroloquinoline quinone (PQQ) for the glucose dehydrogenase. The electrons acquired by the enzyme then may be moved to the electrode by a mediator, which is converted to a reduced form through oxidation of the enzyme. Finally, the reduced form of the mediator, such as the ferrocyanide species of the ferricyanide/ferrocyanide redox pair, is oxidized at the electrode to generate a measurable current.\nThis process may be represented by the following equations:Glucose+EOx==ERed+Product  (1)ERed+nMedOx==nMedRed+EOx  (2)MedRed==MedOx+ne−  (3)where EOx and ERed are the oxidized and reduced forms of the redox center of the enzyme, respectively, while MedOx and MedRed are the oxidized and reduced forms of the mediator, respectively. The product of the enzymatic reaction may be gluconic acid or gluconolactone.\nOne electrochemical method, which has been used to quantify analytes in biological fluids, is coulometry. For example, Heller et al. described the coulometric method for whole blood glucose measurements in U.S. Pat. No. 6,120,676. In coulometry, the analyte (glucose) concentration is quantified by exhaustively oxidizing the analyte within a small volume and integrating the current over the time of oxidation to produce the electrical charge representing the analyte concentration. In other words, coulometry captures the total amount of glucose within the sensor strip.\nAn important aspect of coulometry is that towards the end of the integration curve of charge vs. time, the rate at which the charge changes becomes relatively constant to yield a steady-state condition. This steady-state portion of the coulometric curve forms a relatively flat plateau region in the curve, thus allowing accurate determination of the corresponding current. However, the coulometric method requires the complete conversion of the entire volume of analyte. As a result, this method is time consuming and does not provide the fast results which users of electrochemical devices, such as glucose-monitoring products, demand. Another problem with coulometry is that the small volume of the sensor cell must be controlled in order to provide accurate results, which can be difficult with a mass produced device.\nAnother electrochemical method which has been used to quantify analytes in biological fluids is amperometry. In amperometry, current is measured at the end of a period at a constant potential (voltage) across the working and counter electrodes of the sensor strip. The current is used to quantify the analyte in the biological sample. Amperometry measures the rate at which the electrochemically active species, and thus the analyte, is being oxidized or reduced. Many variations of the amperometric method for biosensors have been described, for example in U.S. Pat. Nos. 5,620,579; 5,653,863; 6,153,069; and 6,413,411. The amperometric method samples the analyte concentration near the electrode surface by measuring the current that is proportional to the diffusion rate and the bulk concentration of the analyte.\nA disadvantage of the amperometric method is the non-steady-state nature of the current after applying a potential. The rate of current change with respect to time is very fast initially and becomes slower as the analysis proceeds due to the changing nature of the underlying diffusion process. Until the consumption rate of the reduced mediator at the electrode surface equals the diffusion rate, a steady-state current cannot be obtained. Thus, measuring a current during a non-steady-state time period may be associated with more inaccuracy than a measurement taken at a steady-state time period.\nOne important aspect of measuring analytes in whole blood samples is the effect of hematocrit. Hematocrit is the volume of red blood cells (RBC) expressed as a percentage of the volume of RBC in a whole blood sample. The hematocrit value for whole blood samples ranges from about 20 to 60% and is typically about 40%.\nReagent biosensors include any system that can detect glucose in a blood specimen via an electrochemical reaction. Examples of reagent biosensors include Ascensia AUTODISC® and Elite® biosensors available from Bayer HealthCare, LLC of Elkhart, Ind.; Precision® biosensors available from Abbott in Abbott Park, Ill.; Accucheck® biosensors available from Roche in Indianapolis, Ind.; and OneTouch Ultra® biosensors available from Lifescan in Milpitas, Calif.\nTypical electrochemical sensor strips contain a working electrode, a counter electrode, and an optional third electrode. A reference potential may be provided to the system by the counter electrode, if configured appropriately, or by the optional third electrode. A reagent layer with an enzyme such as glucose oxidase or glucose dehydrogenase and a mediator such as ferricyanide or ruthenium hexaamine is printed or deposited onto the working electrode or onto the working and counter electrodes with a polymer as the binder.\nExamples of polymers used as the binder of the reagents include CMC (carboxyl methyl cellulose) and PEO (polyethylene oxide). The addition of various types and molecular weights of polymers to the reagent layer may assist in filtering red blood cells, preventing them from coating the electrode surface.\nPreferably, the sensor strip is made by printing electrodes on an insulating substrate using multiple techniques, such as those described in U.S. Pat. Nos. 6,531,040; 5,798,031; and 5,120,420. The reagent can be co-printed onto the working and counter electrodes with a mixture of a glucose oxidizing enzyme such as glucose oxidase, a mediator such as ferricyanide, a hydrophilic polymer such as polyethylene oxide (PEO) and an appropriate buffer, such as a citrate buffer.\nAlternatively, a different reagent chemistry can be either printed or micro-deposited separately onto the working and counter electrodes using the method described in a U.S. provisional patent application filed Oct. 24, 2003, Ser. No. 60/513,817 with the reagent on the working electrode containing the enzyme, the mediator, the polymer and that on the counter electrode containing a soluble redox species, which could be the same as the mediator or different, and a polymer. In one embodiment, the polymer used in micro-deposition is carboxyl methyl cellulose.\nExamples of suitable bench-top electrochemical instruments which may be used for reading reagent biosensors according to the present invention include, but are not limited to, the BAS 100B Analyzer available from BAS Instruments in West Lafayette, Ind.; the CH Instrument Analyzer available from CH Instruments in Austin, Tex.; the Cypress Electrochemical Workstation available from Cypress Systems in Lawrence, Kans.; and the EG&G Electrochemical Instrument available from Princeton Research Instruments in Princeton, N.J. Examples of portable instruments include the Ascensia Breeze® and Elite® meters of Bayer Corporation.\nA biosensor for glucose may have an enzyme and a mediator deposited on the electrodes. The ability of this sensor to measure glucose is affected as the RBC block the diffusion of the relevant reagents within the blood sample. Since the amperometric current is directly proportional to the diffusion of the reduced form of the mediator, the hematocrit will have a significant impact on the accuracy of the glucose measurements. Depending on the hematocrit level in a whole blood sample, the RBC cause a bias in the glucose readings.\nVarious methods and techniques have been proposed in an attempt to reduce the hematocrit effect of the whole blood on the resulting glucose measurements. For example, Ohara et al. in U.S. Pat. No. 6,475,372 disclosed a method of using the ratio of currents from a forward and a reverse potential pulse to compensate for the hematocrit effect in electrochemical glucose measurements. McAleer et al. in U.S. Pat. Nos. 5,708,247 and 5,951,836 disclosed a reagent formulation using silica particles to filter the RBC from the electrode surface, thus reducing the hematocrit effect. Carter et al. in U.S. Pat. No. 5,628,890 disclosed a method of using a wide spacing of the electrodes combined with mesh layers to distribute the blood sample for the hematocrit effect.\nThese conventional techniques for reducing the bias attributable to the hematocrit effect include (a) co-deposition of a polymer to minimize the hematocrit effect, (b) addition of various kinds of fused silica to enforce the filter effect for the polymer layer, (c) compensation coefficients based on the ratio of currents from a forward and a reverse potential pulse, and (d) self-compensation by utilizing the existing solution resistance of the whole blood samples. Although these methods may be useful, conventional glucose sensors continue to exhibit significant analytical bias attributable to the hematocrit effect. Thus, it would be desirable to provide systems for quantifying analytes in biological fluids, in particular the glucose content of whole blood, which reduces bias from the hematocrit effect."}
{"text": "As current molecular target drugs against malignant tumors, there have been developed and used, in practical clinic, many agents such as sorafenib used for liver cancer and renal cancer, cetuximab and bevacizumab used for large intestine cancer, erlotinib and gefitinib used for lung cancer, and trastuzumab used for breast cancer. As target substances of each agent, there can be exemplified cancer kinase, vascular endothelial growth factor (VEGF), epidermal growth factor receptor (EGFR), HER2 protein, and the like. However, in relation to molecular target drugs against malignant tumors, there have also been reported skin disorder, gastrointestinal perforation, interstitial pneumonia, and deaths due to severe side effects such as liver failure. Therefore, appearance of a safer molecular target drug is desired. Meanwhile, although it is said that even healthy persons have thousands of malignant tumor (cancer) cells every day, everyone does not necessarily develop disease. It has been reported that incidence of liver cancer in autoimmune hepatitis (AIH) is about 0.7%/year (Aliment Pharmacol Ther 2006; 24:1197) and is obviously lower than 3%/year in chronic hepatitis C (Ann Intern Med 1999; 131:174).\nIt has been reported that a ribosomal protein RPL29 (Ribosomal protein L29) shows increased expression in cells of various malignant tumors such as large intestine cancer (Non Patent Literature 1), liver cancer (Non Patent Literature 2), gastric cancer (Non Patent Literature 3), thyroid cancer (Non Patent Literature 4) and breast cancer (Non Patent Literature 5). The RPL29 is a membrane protein expressed on a cell surface, and it has been reported that when the expression of the RPL29 is lowered, cells causing apoptosis are increased (Non Patent Literature 6), and cellular differentiation is induced (Non Patent Literature 1). In addition, it has been reported that a ribosomal protein RPS4X (Ribosomal protein S4, X-linked) shows increased expression in cells of various malignant tumors such as large intestine cancer (Non Patent Literature 7) and breast cancer (Non Patent Literature 8). However, it has not been reported that a substance targeting the ribosomal proteins RPL29 and RPS4X improved tumors."}
{"text": "The present invention is directed to street and highway posts generally used to support a signboard and more specifically to a connection adapted to be mounted between a lower post section and an upper post section. The connection is particularly adapted to prevent damage to the lower post section when the post is hit by a vehichle, and in many instances to also prevent damages to the upper post section.\nThe invention is particularly directed to a specific type of connection having a corrugated cross-section and more specifically, to a U-shape channel having two lateral flanges extending outwardly from the U-shaped channel. The connection is provided with a weakened linear section extending partly through the cross-section of the channel so as to allow the tilting of the upper post section relative to the lower post section and a subsequent tearing apart of the connection to maintain the upper post section linked to the lower post section."}
{"text": "This invention relates to an emanator device for delivering an air treatment agent into the surrounding air.\nThe invention relates in particular to an emanator device for delivering an air treatment agent into the surrounding air, for use in an environment which is alternately wet and dry.\nWO 2004/069982 A1 describes an emanator for emanating one or more active component(s), which preferably comprises a surfactant and fragrance. Typically the emanator is for use in a dishwashing machine. The emanator comprises a blister which contains the active component(s) as a concentrated liquid, within a permeable polymeric membrane, through which the active component(s) gradually diffuses. The blister is retained in a cage-like retaining structure. The arrangement, in such a device, is such that the blister is washed over by water in “wet phases” of the dishwashing cycle; and at other times is dispensed into the ambient air. Thus, emanation of the active component(s) occurs without interruption, whether into the ambient air or into the ambient water.\nWO 2005/102141 A1 describes an emanator of a similar type, having an aperture or apertures whose size may be varied, in order to provide different rates of release of the active component(s). As in the device of WO 2004/069982 A1, the emanation of the active component(s) is into the ambient environment, whether into water or into air, and 25 thus proceeds without interruption.\nIt is desirable to fragrance the air in a dishwashing device, to mask the chemical smell sometimes found in a dishwasher, and the odour of dirty dishes and cutlery waiting to be washed. However it is undesirable to add fragrance to the water in the wash cycle, because the crockery and cutlery may have a taint which some users can detect.\nWO 2008/141473 A1 describes a malodour-counteracting device for use in an automatic dishwashing machine. The device addresses a problem, that there is increased evaporation of fragrance during the wash cycle, followed by reduced evaporation between washes. The latter time is often when more, rather than less, fragrance is needed, given that dirty dishes and cutlery are loaded into the dishwashing machine, and may be kept there for some hours, before the dishwashing machine is operated. To solve this problem a device is described having a downwardly-open body holding a gradually-releasable fragrance source, and a larger, upwardly open, dish. The body and the dish overlap, but are not sealed to one another; fragranced air can leave the device though a space between the dish and the body. However, during the wet phase of a dishwasher cycle water collects in the dish and thereby rises over the space, preventing fragranced air from leaving the device. However the fragrance source is not itself immersed in water. Fragranced air can not escape from the device until the level of water in the dish falls, by evaporation, to open the space.\nThe device of WO 2008/141473 A1 achieves the objective of preventing fragrance release during the wet phase, by use of a collected water barrier. However, the device described is aesthetically unappealing, and has functional limitations. The “dish” of the device needs to retain a substantial volume of water in order for the water to close the passage through which fragrance could otherwise escape. To open the passage to its maximum degree appears to require the complete absence of water in the dish. Moreover it is sometimes found that water collected in upwardly-open formations, such as the bases of cups, can contain residues of dirty water, in certain makes of dishwashing machine. This can be expected to occur on occasion, in the upwardly open “dish” of the device of WO 2008/141473 A1. Providing a device which relies upon an upwardly-open dish for alternating closure and non-closure of an emission port or ports is not an optimal technical solution."}
{"text": "Known in the prior art are semiconductor devices which provide volatile storage of information in a dynamic random access memory (RAM) portion and non-volatile storage in a read only memory (ROM) portion. For example, U.S. patent application Ser. No. 654,332, (now U.S. Pat. No. 4,611,309) filed Sept. 24, 1984, on behalf of Patrick Chuang, Robert Yau and Ron Maltiel, entitled \"An Improved Non-Volatile Dynamic RAM Cell\", assigned to the assignee of the instant application discloses one such device.\nOftentimes it is necessary to non-destructively access both the volatile and non-volatile data for such purposes as comparison in image processing and pattern recognition. It is desirable that both of these accesses take place over a short period of time, ideally within one memory clock cycle, because of the large number of data points used in such applications as image processing and pattern recognition. If each such comparison requires several clock cycles, an undesirable time penalty would be imposed during pattern recognition.\nHowever, the memory cells of the prior art are incapable of concurrently accessing both the volatile and non-volatile data without over-writing the volatile data. In addition, in these so-called shadow-memories, the non-volatile portion merely duplicated the data stored in the volatile portion. This duplication was necessary to provide back-up for data which would be lost from the volatile portion during a power failure. Therefore, only one-half of the shadow-memory is being used, the other one-half serving as back-up. On approach overcoming these difficulties is to temporarily store the data from, for example, the volatile portion of the memory during one memory access cycle and then access the data from the non-volatile portion of the memory during the following memory cycle for subsequent comparison. Such approach is described in a co-pending application Ser. No. 717,348, filed Mar. 29, 1985, on behalf of Ron Maltiel and Robert Yau, entitled \"Method and Apparatus for Non-Destructive Access of Volatile and Non-Volatile Data in a Shadow Memory Array\", and assigned to the assignee of the instant application."}
{"text": "Users who want to subscribe to services of a wireless wide area network (“WWAN”) typically have to complete a complicated series of steps. The device (e.g., personal computer, cell phone, or personal digital assistant) through which the services are to be accessed may have a wireless network interface card, such as a Global System Mobile (“GSM”) card. When such a device is purchased by a user, the user typically has not subscribed to services of a service provider, and thus the wireless network interface card is not provisioned to access any service provider. To subscribe to the services of a service provider, the user may have to research the services of various service providers by locating and accessing the web sites of the service providers. Once the user has decided upon a service provider and a service, the user needs to subscribe to the service and obtain a subscriber identity module (“SIM”) card that is provisioned with information provided by the service provider. A SIM card, which contains a subscriber identity that uniquely identifies the SIM card, is provisioned by storing in its memory information relating to registration information and activation service information for the selected service provider, which may be considered a “home” service provider. The registration information and activation service information may be stored on a network server, rather than in the memory of the SIM card. The activation information may include an access point name (“APN”) for a subscribed-to service of the selected service provider. The device uses the registration information to notify a service provider that it is currently on and within range of the service provider's wireless network. The registered device uses the service activation information to notify the service provider that it wants to start accessing a subscribed-to service (e.g., access the internet). The provisioning may also include storing of information for other service providers that may be “home,” “preferred,” or “non-preferred” service providers. The selected service providers may have agreements with the other service providers to provide some of the services, for example, in an area not serviced by the selected service provider. For example, if the services of a home service provider are not currently available to a device, the device may try to activate a service of a preferred or non-preferred service provider. The user can then install the SIM card onto the GSM card to access the services of the service provider. However, to obtain such a SIM card, the user may need to go to a retail outlet of the service provider and subscribe to the service. Alternatively, the user may obtain the SIM card via mail order. Regardless of how the SIM card is obtained, the user typically needs a SIM card that is provisioned by a service provider in accordance with the selected service.\nSome vendors of devices may provide a purchase option by which a purchaser of a device can subscribe to a service of a wireless service provider at the time the device is purchased. The purchaser may select a service provider and service that is made available by the vendor. The vendor coordinates the subscribing of the purchaser to the service and the provisioning of the SIM card with service provider and service information. The vendor then installs the provisioned SIM card into the device before it is delivered to the purchaser.\nThe selection of a service provider and a service at the time of purchase of the device limits both the timing of the subscription and the choice of service providers and services. Some vendors may require a purchaser to subscribe to a service at the time of purchase and limit the purchaser's options to certain service providers and services. At the time of purchase, a purchaser may not be ready to subscribe to a service (but may want the option of doing so later) or may not want to subscribe with the service providers or services that the vendor makes available. If the purchaser does not subscribe at the time of purchase, then later subscription requires the complicated series of steps as described above. It would be desirable to have a less complicated way of subscribing to wireless services that would offer a range of service providers and services and not necessarily be limited to subscribing at the time of purchase of a device."}
{"text": "This section provides background information related to the present disclosure and is not necessarily prior art.\nIn an attempt to reduce the quantity of NOx and particulate matter emitted to the atmosphere during internal combustion engine operation, a number of exhaust aftertreatment devices have been developed. A need for exhaust aftertreatment systems particularly arises when diesel combustion processes are implemented. Typical aftertreatment systems for diesel engine exhaust may include a hydrocarbon (HC) injector, a diesel oxidation catalyst (DOC), a diesel particulate filter (DPF), and a selective catalytic reduction (SCR) system.\nSCR technology has been used in conjunction with reducing nitrogen oxides present in the exhaust of combustion engines. Some of these systems are constructed using urea-based technology including a container for storing a reductant (e.g., urea) and a delivery system (including an injector, for example) for transmitting the reductant from the container to the exhaust stream. A mixer is typically provided for mixing the injected reductant with the exhaust gas before the reductant reaches a catalyst with which the reductant reacts. While these systems may have performed well in the past, it may be desirable to provide an improved mixer to more efficiently and effectively mix the reductant with the exhaust stream."}
{"text": "This invention relates to an evolutionary controlling system, and particularly to that for controlling characteristics of a subject in an evolutionary manner.\nHeretofore, when a control system or control characteristics of a subject, such as vehicles and electrical appliances, is designed, imaginary users are selected, and the users' preferences and their using conditions are taken into consideration. The characteristics of the subject are determined in such a way as to adapt the subject to users in as broad a range as possible.\nHowever, each individual user has a particular and unique personality, and thus, their preferences are diverse. Thus, there is a problem in that even if imaginary users are selected to develop and design a product for users by presuming the users' preference, it is impossible to satisfy all of the users of the product.\nIn order to solve the above problem, prior to purchase of a product, a prospective user is requested to determine whether or not the product is satisfactory to the user after checking the characteristics of the product in light of the user's preferences. However, it is troublesome for the user to check the characteristics of the product before purchase. Further, because a series of products are often operated or controlled by characteristics common in the products, although the design of the product is changed depending on the user's preferences, the user may not like other operational characteristics. Thus, although the design is appealing to some prospective users, the users may not purchase the product since the operational characteristics do not appeal to them. In other words, there is another problem in that the range of users is limited and depends on the operational characteristics.\nAn objective of the present invention is to provide an evolutionary control system to construct characteristics which can satisfy plural users."}
{"text": "For displaying graphic images on a cathode ray tube display monitor, personal computer systems generally include a video adapter circuit board housed within the computer system case. The video adapter circuit board typically includes a bus interface circuit, a dedicated processor, dual-port random access memory (VRAM), and a digital-to-analog converter. The video adapter circuit board is typically controlled by the central processing unit (CPU) of the computer system and receives digital data representative of an image to be displayed from the computer system memory or from an input device coupled to the computer system bus.\nTypically, video adapter circuit boards generate red, green and blue (RGB) analog component video signals in addition to horizontal (HSYNC) and vertical (VSYNC) synchronizing signals. These signals are generally formatted according to a selected one of several commonly utilized computer scanning formats. These formats determine the resolution of the video image and include video graphics array (VGA), super video graphics array (SVGA), color graphics adapter (CGA), extended graphics adapter (EGA), and extended graphics array (XGA) formats. A cathode ray tube (CRT) adapted to the selected scanning format receives the analog RGB component video signals and synchronizing signals and provides a display of the image.\nDue to concerted efforts of display manufacturers, flat panel display monitors are becoming increasingly available. Flat panel display monitors offer several advantages over CRT display monitors. For example, flat panel display monitors are typically lighter in weight, occupy less space, consume less power, emit less radiation and have less glare than do their CRT counterparts. Commercially available flat panel display monitors are commonly liquid crystal displays that utilize active-matrix addressing. These include thin film transistor (TFT) and supertwisted-nematic (STN) flat panel display monitors.\nTFT and STN flat panel display monitors operate on markedly different principles than do CRT display monitors and, therefore, require distinctly different driving signals. TFT and STN flat panel display monitors utilize the optical properties of liquid crystal which are different when measured in a direction parallel to the axis of the crystals in comparison to when measured in a direction perpendicular to the axis of the crystals. Liquid crystal is sandwiched between two polarizing filters having polarizing directions oriented 90 degrees apart. The inside surface of the display adjacent to the liquid crystal is treated to cause the liquid crystals to align such that they appear transparent. When an electric field is applied to an area of the display, liquid crystals in that area are reoriented by the electric field such that they appear opaque. By adding a backlight and color filters, a color image can be displayed. In active matrix displays, a transistor is dedicated to control the electric field for each pixel. TFT displays are characterized in that under an applied electric field, the crystals are reoriented by twisting up to 90 degrees. In an STN display, the crystals are reoriented by twisting up to 270 degrees under the applied electric field. STN displays generally have higher contrast and greater viewing angles that TFT displays.\nBecause a CRT accepts RGB signals, whereas a flat panel display accepts signals utilizing a digital protocol, these devices cannot merely be interchanged in a computer system. Rather, the computer must typically be disassembled and appropriate video card drivers are exchanged. It is anticipated, however, that a user of an existing personal computer system will wish to utilize a flat panel display monitor with the personal computer system that is configured for a CRT. It is further anticipated that such a user will be discouraged from so doing if hardware and/or software modifications to the personal computer are required or if considerable technical knowledge is required on the part of the user.\nTherefore, what is needed is a circuit for adapting a personal computer system for driving a flat panel display monitor without requiring modification of the personal computer system hardware or software. Further, what is needed is a circuit that can determine the scanning format of analog RGB signals generated by a personal computer and, based upon the determination, convert the analog RGB signals into signals appropriate to driving a flat panel display monitor. What is also needed is an ability to adjust display parameters for the flat panel display monitor, such as brightness and contrast, without requiring modification of the personal computer system."}
{"text": "1. Field of the Invention\nThe present invention relates to a folding cup holder attached to an interior of a passenger compartment in an automobile, and more particularly to a folding cup holder having a structure which does not easily break even when accidental external force is applied, and is easy to be restored.\n2. Description of the Related Art\nVarious kinds of equipment such as a cup holder are mounted to the interior of the passenger compartment in an automobile for the convenience of the passenger. There are cup holders which are built into the interior, however, there is proposed cup holders of an add-on type. The add-on type cup holder is proposed as a folding structure, which does not obstruct the passenger when it is not in use.\nIn the folding cup holder, it is necessary to make an arm for holding a cup and a structure for supporting the arm narrow and small, for the purpose of achieving both of the matters of the cup holder being compactly foldable and the cup holder being expandable largely. Accordingly, it is necessary to employ a countermeasure to prevent the cup holder from being broken in the case where accidental external force is applied. Japanese Patent Application Laid-open No. 2000-8687 shows an art of the folding cup holder. This cup holder is fixed to a wall surface or the like of the interior, can be folded compactly when not being used, and can be expanded forward when being used. The structure is such that a structure for supporting a bottom plate is cancelled so as to prevent the structure from being broken when accidental external force is applied.\nAccording to the proposed folding cup holder mentioned above, when the structure for supporting the bottom plate is cancelled due to the accidental external force, the bottom plate drops down. In order to restore this, it is necessary to reassemble the bottom plate in the supporting structure again according to a troublesome operation. It is hard to smoothly execute this operation without considerable practice, and in particular, it is even harder to execute while driving the vehicle.\nThe present invention is made in view of the above problem, and an object of the present invention is to provide a folding cup holder having a structure which does not easily break even when accidental external force is applied and can be easily restored.\nAccording to the present invention, a folding cup holder is provided with a base plate including one or more through holes; a holder arm for holding a cup, the holder arm being pivotally supported to the base plate; and a tray including one or more latch arms respectively inserted to the through holes, and axially supported to the base plate. The latch arms and the through holes engage with each other to hold the tray in a position for holding a cup and an engagement thereof is easily cancelled in a case that excessive force is applied to the tray.\nAccording to the structure mentioned above, when excessive force is applied to the base plate, the latch arms easily displace from the through holes and the base plate becomes free to swing, so that it is possible to prevent the base plate and a bearing thereof from being broken. Further, at a time of making the structure in an initial state, the latch arms easily engage with the through holes simply moving the base plate upward by hand and inserting the latch arms to the through holes. Accordingly, it is possible to easily restore to the initial state without a complex operation.\nPreferably, each of the through holes is provided with a leaf spring for catching the latch arm and the leaf spring is fitted to an edge thereof. The engagement and the cancellation of the latch arms mentioned above can be more smoothly executed by engaging the latch arms with the leaf springs in place of the through holes. Accordingly, as well as a risk of breakage being reduced, the restoring operation can be more easily executed.\nMore preferably, a distal end of each of the latch arms is formed in a fork shape. Since the fork shape increases flexibility of the distal end of the latch arm, the risk of breakage is reduced, and the restoring operation can be more easily executed.\nFurther preferably, through holes and the latch arms are constituted in a pair respectively and proximal ends of the latch arms are integrally connected by a rib which has a U-shaped end portion. The latch arms are reinforced with the rib so that the latch arms are hard to be broken."}
{"text": "1. Field of the Invention\nThe present invention relates to improvements of a surface-light-emitting device including a light-generating layer and a light resonator constituted by two multi-film reflecting layers between which the light-emitting layer is interposed.\n2. Discussion of Related Art\nThere is known a surface-light-emitting device including a semiconductor substrate and a plurality of semiconductor layers which are grown on the semiconductor substrate. The semiconductor layers comprise a light-generating layer (active portion), and a light resonator constituted by two multi-film reflecting layers (DBR: distributed-Bragg reflectors) which are located on the opposite sides of the light-generating layer, to reflect a light generated by the light-generating layer. The semiconductor layers further comprise an outermost semiconductor layer which is located on one of the opposite sides of a laminar semiconductor structure consisting of the above-indicated semiconductor layers, which one side is remote from the substrate. The outermost semiconductor layer has a light-emitting surface from which the light generated by the light-generating layer and resonated by the light resonator is emitted. Examples of such a surface-light-emitting device include a resonant-cavity light-emitting diode and a surface-light-emitting laser.\nJP-A-2000-299492 discloses a light-emitting diode of quantum-well type, as an example of the surface-light-emitting device. This light-emitting diode includes a plurality of semiconductor layers laminated on a substrate. The semiconductor layers comprise a light-generating layer of a quantum-well structure having a thickness not larger than the wavelength (100 xc3x85, i.e., 10 nm) of the electron wave, and a light resonator consisting of two reflecting layers which are located on the respective opposite sides of the light-generating layer and which reflect a light generated by the light-generating layer. By applying an electric current between a pair of electrodes formed on the respective opposite surfaces of a laminar semiconductor structure consisting of the semiconductor layers, the light generated by the light-generating layer is emitted from one of the opposite surfaces of the laminar semiconductor structure which is remote from the substrate. In the light-emitting diode disclosed in the above-identified publication, an electron wave within the light-generating layer and an optical wave within the light resonator are coupled together, so that the light-generating layer generates a light only in a resonance mode. This phenomenon so-called xe2x80x9ca cavity QED effectxe2x80x9d permits the emitted light to have a high degree of directivity and a narrow line width, and prevents total reflection of the light by the crystal surface, assuring an advantage of a high degree of external quantum efficiency.\nA known surface-light-emitting device as described above includes, for instance, an n-GaAs semiconductor substrate, and a first multi-film reflecting layer, a light-generating layer and a second multi-film reflecting layer which are successively grown on the substrate. The first multi-film reflecting layer consists of a multiplicity of n-(Al)GaAs/Al(Ga)As films superposed on each other such that the (Al)GaAs films and the Al(Ga)As films are alternately arranged. The ratios of the elements indicated in the parentheses are suitably selected, and may be zero. The second multi-film reflecting layer consists of a multiplicity of p-(Al)GaAs/Al(Ga)As films superposed on each other such that the (Al)GaAs films and the Al(Ga)As films are alternately arranged. The parentheses have the same meaning as described just above. Where the wavelength of the light to be emitted from the surface-light-emitting device is as short as about six hundred and some tens of nanometers (nm), the concentrations of Al of the (Al)GaAs films must be made relatively high in order to reduce the light absorptance. However, an increase in the Al concentration undesirably reduces or substantially zeros the refractive index at the boundaries of the alternately arranged films, leading to difficulty in obtaining a sufficiently high degree of reflectivity of the multi-film reflecting layers. Another problem encountered in the known surface-light-emitting device relates to a change of a gas used to form the first multi-film reflecting layer (formed on the side of the substrate) and the light-generating layer (active portion of the device). Described in detail, epitaxial growth of the AlGaInP/GaInP light-generating layer on the first multi-film reflecting layer is initiated by changing the gas from arsine AsH3 which has been used for growing the As-based material of the first multi-film reflecting layer, to phosphine PH3 to be used for the epitaxial growth of the light-generating layer on the first reflecting layer. However, this arsine-to-phosphine change of the gas takes a relatively long time, causing an increased surface roughness of the first reflecting layer, and oxidization of Al, making it difficult to form a high-quality active portion on the first reflecting layer, that is, the light-generating layer having a high quality on the first reflecting layer.\nOn the other hand, it has been proposed to form the first and second (n- and p-) multi-film reflecting layers consisting of (Al)GaInP/Al(Ga)InP films. The reflecting layers according to this proposal exhibit a high degree of reflectivity even where the light to be emitted has a considerably short wavelength. In addition, the epitaxial growth of the AlGaInP/GaInP light-generating layer on the first reflecting layer does not require a change from arsine AsH3 to phosphine PH3, before the epitaxial growth of the light-generating layer on the first reflecting layer. Accordingly, the light-generating layer serving as the active portion of the surface-light-emitting device can be advantageously given a high degree of crystallinity. However, the use of the p-(Al)GaInP/Al(Ga)InP second multi-film reflecting layer gives rise to a relatively large degree of carrier blockage, and an accordingly high electric resistance, and results in reduced enclosure of electrons, so that the operating efficiency of the surface-light-emitting device is deteriorated where the device is designed to emit a light having a relatively short wavelength. Another problem arises from the use of the (Al)GaInP/Al(Ga)InP films for the first and second (n- and p-) multi-film reflecting layers, where a mesa structure is formed by selective etching effected on the laminar semiconductor structure grown on the semiconductor substrate. That is, it is extremely difficult to practice a wet-etching process to form the mesa structure. Although the application of the wet-etching process is not impossible, the yield ratio in the production of the surface-light-emitting device is considerably low where the wet-etching process is used to form the mesa structure. In addition, the difficulty to form the mesa structure by the wet-etching process leads to difficulty to form a current-blocking region by subjecting the peripheral portion of the mesa structure to an oxidizing treatment.\nThe present invention was made in view of the prior art described above. It is therefore a principal object of the present invention to provide a surface-light-emitting device which includes multi-film reflecting layers having a high degree of reflectivity with respect to a relatively short wavelength of light, and a light-generating layer having a high degree of crystallinity, and which has a low electric resistance and exhibits a high operating efficiency. It is an optional object of this invention to provide a surface-light-emitting device including a laminar semiconductor structure which is formed on a semiconductor substrate and which has a mesa structure formed by wet-etching and having a current-blocking region formed by an oxidizing treatment.\nThe principal object indicated above may be achieved according to the principle of the present invention, which provides a surface-light-emitting device including a semiconductor substrate, and a laminar semiconductor structure consisting of a plurality of semiconductor layers formed by epitaxial growth on the semiconductor substrate, the laminar semiconductor structure including a light-generating layer, and two multi-film reflecting layers between which the light-generating layer is interposed and which cooperate to constitute a light resonator for reflecting a light generated by the light-generating layer, the laminar semiconductor structure having a light-emitting surface at one of opposite ends thereof remote from the semiconductor substrate, the light generated by the light-generating layer being emitted from the light-emitting surface, characterized in that the two multi-film reflecting layers consist of a first multi-film reflecting layer formed principally of AlGaInP on the semiconductor substrate, and a second multi-film reflecting layer formed principally of AlGaAs on one of opposite sides of the light-generating layer which is remote from the semiconductor substrate.\nIn the present surface-light-emitting device, the first multi-film reflecting layer formed on the semiconductor substrate, which is required to have a comparatively high degree of reflectivity, is formed principally or essentially of AlGaInP, while the second multi-film reflecting layer located on the side of the light-emitting surface is formed principally or essentially of AlGaAs. The first multi-film reflecting layer formed principally of AlGaInP has a large difference between refractive index values of adjacent ones of the alternately arranged films, and a low degree of light absorptance, so that the first multi-film reflecting layer exhibits a high reflectivity, even where the light to be emitted has a relatively short wavelength. In addition, the formation of the light-generating layer (formed principally of AlGaInP/GaInP, for example) by epitaxial growth on the first multi-film reflecting layer does not require a change of a gas to form the light-generating layer after the first reflecting layer, so that the light-generating layer is given a high degree of crystallinity. Further, the first multi-film reflecting layer formed principally of AlGaInP has a lower electric resistance than a reflecting layer formed principally of AlGaAs, and the light-generating layer is interposed between the first reflecting layer formed principally of AlGaInP and the second reflecting layer formed principally of AlGaAs on the side of the light-emitting layer, so that a carrier-enclosing region for enclosing the carrier can be suitably formed upstream of the barriers as seen in the direction of movement of the carrier, assuring a high degree of operating efficiency of the device.\nAccording to one preferred form of this invention, the semiconductor substrate is formed of n-GaAs, and the first multi-film reflecting layer consists of n-(Al)GaInP films and n-Al(Ga)InP films which are alternately laminated on each other, while the second multi-film reflecting layer consists of p-(Al)GaAs films and p-Al(Ga)As films which are alternately laminated on each other. In this case, the light-generating layer is formed of AlGaInP/GaInP. In the device according to this preferred form of the present invention, the first multi-film reflecting layer consisting of the alternately arranged n-(Al)GaInP films and n-Al(Ga)InP films has a high refractive index and a low degree of light absorptance, and exhibits a high degree of reflectivity even where the light to be emitted has a comparatively short wavelength, and the formation of the light-generating layer of AlGaInP/GaInP by epitaxial growth on the first multi-film reflecting layer of n-(Al)GaInP/Al(Ga)InP does not require a change of the gas to form the light-generating layer on the first reflecting layer, so that the light-generating layer is given a high degree of crystallinity. Further, the first multi-film reflecting layer consisting of the alternately arranged n-(Al)GaInP films and n-Al(Ga)InP films has a lower carrier blockage and a lower electric resistance than a reflecting layer consisting of alternately arranged n-(Al)GaAs films and n-Al(Ga)As films, and the light-generating layer is interposed between the first reflecting layer consisting of the alternately arranged n-(Al)GaInP films and n-Al(Ga)InP films and the second reflecting layer consisting of the alternately arranged p-(Al)GaAs films and p-Al(Ga)As films, so that carrier-enclosing regions Ke and Kh for enclosing the carrier can be suitably formed upstream of the barriers as seen in the direction of movement of the carrier (electrons and holes), assuring a high degree of light-emitting efficiency of the device. Where the laminar semiconductor structure grown on the semiconductor substrate includes a mesa structure formed by etching an end portion of the structure, as described below in detail, an ordinary etching process may be used for easy formation of the mesa structure, since the second multi-film reflecting layer formed of p-(Al)GaAs/Al(Ga)As does not include phosphorus (P). Further, a non-electrically-conductive current-blocking region can be easily formed by oxidizing aluminum (Al) in an exposed portion of the mesa structure which includes the second reflecting layer formed of p-(Al)GaAs/Al(Ga)As.\nAccording to another preferred form of this invention, the laminar semiconductor structure includes a mesa structure which has the light-emitting surface and the second multi-film reflecting layer, the mesa structure being formed by etching an end portion of the laminar semiconductor structure, which end portion is laminated on the light-generating layer, the mesa structure having a smaller cross sectional area taken in a plane parallel to the light-emitting surface, than the other portion of the laminar structure, the second multi-film reflecting layer comprising a generally annular non-electrically-conductive current-blocking region which is formed by oxidizing an exposed portion of the mesa structure which is a peripheral part as seen in the above-indicated plane. In this arrangement, the generally annular non-electrically-conductive current-blocking region can be easily formed by oxidizing the peripheral portion of the mesa structure, since the peripheral portion of the mesa structure is exposed in the radial direction of the mesa structure.\nAccording to a further preferred form of this invention, the light-generating layer includes at least one active layer of quantum-well structure each interposed between adjacent two barrier layers, each active layer and the adjacent two barrier layers being located between the first and second multi-film reflecting layers such that each active layer is aligned with an antinode of the light resonated by the light resonator. In the present arrangement, an electron wave within the light-generating layer and an optical wave within the light resonator are coupled together, so as to provide a so-called xe2x80x9ccavity QED effectxe2x80x9d permitting the light-generating layer to generate a light only in a resonance mode. Accordingly, the light emitted from the light-emitting surface has a high degree of directivity and a narrow line width, without total reflection of the light by the monocrystalline semiconductor surface, assuring a high degree of quantum efficiency and a large optical output."}
{"text": "The present invention relates generally to a storage system for managing storage apparatuses and, more particularly, to a method and a system for deploying a virtual hard disk file to a virtual volume known as thin provisioning.\nThe use of virtual hard disk (VHD) has been widespread in accordance with the growth of server virtualization technology. In addition, the virtual appliance which packs virtual server specification and the VHD into one package has also grown in usage. The VHD can be recognized as a logical volume by the virtual machine even though the VHD is recognized as just a file by the physical server because the hypervisor converts the VHD file dynamically and presents it as a volume to the virtual machine. The VHD has a thin provisioning functionality that grows its actual capacity on demand. On the other hand, the thin provisioning functionality has also been implemented in a storage apparatus such as the Hitachi Dynamic Provisioning of Hitachi USP/V. In that case, the virtual volume is presented as a logical unit that can be accessed through a fiber channel mainly from the physical server.\nMigration between the virtual server and physical server is increasingly popular based on the requirement for performance and availability of the application system. However, the physical server cannot access the VHD as a logical drive. Conversion of the VHD to a logical device is necessary for the use of the VHD by the physical server. In that case, if the VHD is converted to a fixed-size logical device, it creates a problem that the capacity efficiency decreases. Furthermore, in the case where the VHD is converted to a virtual device (logical device with thin provisioning), if there is a gap between the page size of the virtual device and the block size of the VHD, the capacity efficiency is not fully realized."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display (LCD), and more particularly, to an LCD and its fabrication method with enhanced reliability.\n2. Discussion of the Related Art\nGenerally, an LCD has a lower substrate having a thin film transistor (TFT) and a pixel electrode, an upper substrate having a color filter for displaying color and a common electrode, and a liquid crystal filled between the lower and upper substrates. Here, glass or crystalline is generally used as the material for the substrates.\nA crystal substrate facilitates the use of a high temperature process, but is rarely used to fabricate a large-sized display. Also, a crystal substrate is too expensive.\nA glass substrate, however, is relatively inexpensive, but cannot be used in a high temperature process, and has a reduced reliability from stress. That is, the stress on the glass substrate is different when depositing a chrome layer on the glass substrate as opposed to depositing an amorphous silicon layer on the glass substrate.\nFIGS. 1A and 1B illustrate a stress characteristics when chrome is generally deposited on a glass substrate. FIGS. 2A and 2B illustrate stress characteristics when amorphous silicon is generally deposited on the glass substrate.\nReferring to FIG. 1A, when a chrome layer 2 is deposited on the glass substrate 1, the substrate takes on a tensile stress. Thus, the glass substrate 1 is bent upward, as shown in FIG. 1B.\nReferring to FIG. 2A, when an amorphous silicon layer 3 is deposited on the glass substrate 1, the glass substrate 1 is bent downward, as shown in FIG. 2B.\nAccordingly, when a data line is formed of chrome in the LCD, the glass substrate is bent by the tensile stress as shown in FIG. 1B, but tends to return to its original state. Hence, a weak portion in the data line may become open.\nThe conventional LCD is described below with reference to the attached drawings.\nGenerally, an upper substrate of the LCD includes an active display area where a TFT and pixel electrode are arranged, a gate line pad for applying a driving signal to a TFT gate electrode of the active display area, and a data line pad for applying a data signal to a TFT source electrode. The TFT and pixel electrode of the active display area is illustrated in FIG. 3. In particular, the TFT is formed at the intersection of the gate line and data line, where the gate line is connected to a gate electrode, the data line is connected to a source electrode, and the drain electrode is formed to be connected to an ITO pixel electrode.\nHere, for the active area of the TFT, the amorphous silicon layer is deposited between the gate electrode and source/drain electrode, and the active area is formed under the data line and is wider than the data line.\nA method for forming the conventional LCD is described below with reference to FIGS. 4A-4E, which are cross-sectional views taken along the line I--I' in FIG. 3.\nAs illustrated in FIG. 4A, a metal such as aluminum (Al) is deposited on a glass substrate 1, and then selectively etched to form a gate line including a gate electrode 4. As illustrated in FIG. 4B, a gate insulating layer 6 (such as a silicon nitride layer), an amorphous silicon layer 7, and a high concentration n-type amorphous silicon layer 8 are sequentially deposited on the overall surface of the glass substrate 1 including the gate line. As illustrated in FIG. 4C, the amorphous silicon layer 7 and high concentration n-type amorphous silicon layer 8 are selectively removed except on the active area and the data line region of TFT.\nAs illustrated in FIG. 4D, a chrome layer 9 is deposited on the overall surface of the substrate, and then selectively removed except on the high concentration n-type amorphous silicon layer 8 of the data line region and on both sides of the amorphous silicon layer 7 and high concentration n-type amorphous silicon layer 8 of the active area to form the data line and source/drain electrode of the TFT. The exposed center portion of the high concentration n-type amorphous silicon layer 8 is removed using the chrome layer of the source/drain electrode as a mask. Here, the data line and source electrode are integrally formed (not shown).\nAs illustrated in FIG. 4E, a passivation layer 10 is formed on the overall surface of the substrate having the chrome layer 9. Then, the passivation layer on the chrome layer 9 of the drain electrode is selectively removed to form a contact hole. A transparent conductive layer is deposited on the overall surface and then selectively etched to remain only on the pixel area to form a transparent electrode 11. The transparent electrode 11 of the pixel area is electrically connected to the chrome layer 9 of the drain electrode through the contact hole.\nHowever, the conventional LCD has problems as follows.\nAs the chrome layer is deposited in a state where the amorphous silicon layer and high concentration n-type amorphous silicon layer are formed on the active area and data line area, (with the remaining portion removed) the glass substrate takes on a tensile stress and becomes bent. Moreover, because the glass substrate also has a tendency to return to its original (unbent) state, the data line becomes open. Therefore, the reliability of the LCD is decreased."}
{"text": "The present invention relates to an endoscope reprocessor for the automated reprocessing of an endoscope. In particular, the present invention relates to an endoscope reprocessor having a rack for holding an endoscope. The invention also relates to an endoscope reprocessor which minimizes the degradation and consumption of disinfectant and is able to shorten the reprocessing time.\nIn the present specification, “reprocessing” an endoscope refers to performing a cleaning and disinfecting process on an endoscope, which process includes specifically a step for cleaning an endoscope with a liquid detergent (cleaning step), a step for disinfecting with a liquid disinfectant the endoscope that has been cleaned with detergent (disinfecting step), and a step for rinsing the endoscope that has been disinfected with disinfectant (rinsing step). The term “cleaning” refers to the cleansing or washing of an endoscope with a liquid detergent that is carried out in the cleaning step.\nAs is commonly known, the endoscope is an instrument used to insert a tube into the human body or another living body and carry out diagnosis, sample collection and treatment within the body, particularly on internal organs.\nThe endoscope is basically composed of an insertion tube which is inserted into the body, a control unit for controlling the insertion tube and for controlling air and water supply, etc. in the endoscope, a connector which connects to an air supply and suction pump (light guide (LG) connector), and a universal cord (LG soft portion) which connects the connector with the control unit and the insertion tube.\nEndoscopes are generally not disposable, single-use devices; rather, they are used repeatedly and on more than one patient. For an endoscope to be repeatedly used in this way, it must be meticulously reprocessed after each use in order to prevent endoscopically transmitted bacterial infections. To this end, various types of endoscope reprocessors for the automated reprocessing of endoscopes have been developed for commercial use.\nSuch apparatuses for reprocessing endoscopes include, for example, the endoscope reprocessor described in JP 2003-135396 A.\nJP 2003-135396 A describes an apparatus for reprocessing an endoscope inside a basin, which apparatus is characterized by having a tray which is insertably disposed at the interior of the basin and on which the endoscope is set. The reprocessor of JP 2003-135396 A is used to reprocess an endoscope by setting the endoscope on the tray outside of the basin, inserting the tray into the basin, and circulating liquid detergent and liquid disinfectant inside the basin.\nAnother example of a known endoscope reprocessor is the apparatus disclosed in JP 2006-68095 A wherein, after the insertion tube and the universal cord of the endoscope have been coiled and the endoscope has been placed in a basin, the endoscope is reprocessed by circulating liquid detergent and liquid disinfectant over a pathway that includes the basin (or by immersion only in the case of the disinfectant).\nSuch an endoscope reprocessor generally reprocesses an endoscope by carrying out a cleaning step in which the endoscope is cleaned with a liquid detergent and may also be rinsed with water, a disinfecting step in which the endoscope is disinfected with a liquid disinfectant, and a rinsing step in which the endoscope is rinsed with water such as tap water to remove disinfectant and the like.\nJP 10-305011 A discloses an endoscope reprocessor which reprocesses only the insertion tube of the endoscope. The insertion tube is placed in a tubular main body of the reprocessor, following which a cleaning step that involves rinsing with water and the use of a liquid detergent and a disinfecting step that involves the use of a liquid disinfectant are carried out, after which drying is carried out by the passage of pressurized air.\nIn the respective steps carried out in an endoscope reprocessor, i.e., the cleaning step, disinfecting step and rinsing step, reprocessing is carried out by circulating the reprocessing liquids not only over the exterior of the endoscope, but also inside the channels at the interior of the endoscope, such as the forceps channel and the air/water channel.\nOwing to the difficulty of discharging liquids from these channels inside the endoscope, after reprocessing is finished, the reprocessing liquids are forced out by introducing air into the channels. Once reprocessing is completed, reprocessing liquids remaining inside the basin which houses the endoscope are drained from the basin gravitationally or with a drain pump."}
{"text": "1. Field of the Invention\nThe present invention relates to a lace buckle device with innovatory convenience for firmly fastening or releasing two ends of a lace, and more particularly to a lace buckle device which can effectively prevent two ends of a lace from being pulled out by force.\n2. Description of the Related Art\nA rope or a lace is commonly used in our daily life. A rope can be used to tie up the mouth of a bag. A lace is applied to tie the mouth of a belted topcoat or a lace shoe for a wearer to wear comfortably. Generally, two ends of a rope or a lace would be tied to form a slipknot so as to secure the rope or lace in position. When one end of the rope or lace being pulled, the slipknot will easily be untied. However, there are many drawbacks by means of tying a slipknot for securing a rope or lace in position. For example, it wastes time and is troublesome to tie a slipknot, and a slipknot is easy to be loosened or to become a fast knot and is not beautiful.\nA conventional lace buckle device has been invented to overcome the drawbacks of tying a slipknot to fasten two ends of a rope or a lace. Referring to FIG. 8, the conventional lace buckle device includes a buckle base 80 and an operative rod 90. An opening 81 is formed on the top of the buckle base 80. A central hollow 82 is provided in the interior of the buckle base 80, but the central hollow 82 does not penetrate through the bottom of the buckle base 80. Two opposite passing holes 83 are provided in the side wall of the buckle base 80. A top tray 91 is provided on the top of the operative rod 90, and a protruding rod 92 is provided on the bottom of the top tray 91. A passing hole 921 is provided to penetrate through the front and rear side walls of the protruding rod 92, and below each end of the passing hole 921 is provided with a stopping strip 922. A protruding block 923 is provided on the bottom of the protruding rod 92.\nReferring to FIG. 9, while being assembled, a spring 93 is placed into the central hollow 82 of the buckle base 80, and then the operative rod 90 is inserted in the central hollow 82 with the direction of the passing hole 921 being in accordance with that of each passing hole 83. At the same time, each stopping strip 922 is forced to pass the opening 81 and stop against each passing hole 83 so as to make the protruding rod 92 be caught in the central hollow 82. Because of the elasticity of the spring 93, the passing holes 83 and the passing hole 921 are slightly staggered, and the passing holes 83 do not respectively wholly lap over each end of the passing hole 921. After the top tray 91 of the operative rod 90 being pressed down, the operative rod 90 will move downward to make the passing hole 921 accord with the passing holes 83, and then two ends of a rope or a lace 94 can penetrate through the overlapped passing holes 83 and 921. When the top tray 91 being released, referring to FIG. 10, two ends of a rope or a lace 94 will be fastened to prevent from being pulled out by means of the stagger of the passing holes 83 and 921.\nHowever, there are still drawbacks in the above-mentioned conventional lace buckle device although it is more convenient to be operated and beautiful to look than a slipknot tied by hands. The conventional lace buckle device can function to fasten or release two ends of a lace by means of the spring 93 and the operative rod 90 which cooperate to change the positions of the passing hole 921, but, when two ends of the lace being clamped by the spring 93 and the operative rod 90, they can still be easily pulled out of the buckle device by forceful pulling. It can not have the perfect efficiency to fasten two ends of a lace.\nTherefore, the present invention is to provide a lace buckle device that substantially obviates the drawbacks of the related conventional art.\nAn object of the present invention is to provide a lace buckle device that is convenient for operating to fasten or release two ends of a lace.\nAnother object of the present invention is to provide a lace buckle device whose catching rims of both slant holes improve perfect efficiency to prevent two ends of a lace from being pulled out by forceful pulling while in clamped condition.\nTo achieve these advantages, a lace buckle device in the present invention includes a buckle base, an operative rod and a spring. The buckle base is provided with a central vertical hollow, and a lateral hole is provided to penetrate through the centers of the front and rear walls of the buckle base. Two slant holes are respectively provided at the left and right sides of the lateral hole and respectively slantwise extend to and through the left and right side walls of the buckle base. A catching rim is formed on each shorter section of each slant hole in the left side wall or right side wall. The operative rod is provided with a lateral passing hole penetrating through the front and rear walls. After the spring and the operative rod being successively inserted in the central vertical hollow, each end of a lace can be inserted into each end of the slant hole in the left side wall or right side wall and out of each other end of the slant hole in the front wall, and then both ends of a lace are reversely inserted into the front end of the lateral hole and penetrate through the lateral passing hole of the operative rod and out of the rear end of the lateral hole of the buckle base so as to make two ends of a lace be firmly fastened by the clamp of the spring and the operative rod and the catch of the catching rims."}
{"text": "Traditional women's undergarments, namely brassieres, generally have a front portion with supportive cups, wings extending from the front portion surrounding the wearer's ribcage, a hook and eye closure to connect the two wings, and two shoulder straps. While traditional bras are practical for everyday use under daily wear, such as t-shirt, blouses, sweaters, etc., these traditional bras are not suitable for all purposes, e.g., for wear with backless garments such as formal gowns.\nOne solution has been a strapless and backless bra as disclosed in U.S. Pat. No. 6,397,391. While there are no wings of the bra to be exposed under a backless garment, this configuration provides little support to the breasts, resulting in an undesirable appearance. There have been other attempts to solve this problem but with little success."}
{"text": "Aluminum metal has been produced for 90 years in the Hall cell by electrolysis of alumina in a molten cryolite salt electrolyte bath operating at temperatures in the range of 900.degree.-1000.degree. C. The reactivity of the molten cryolite, the need for excellent electrical conductivity, and cost considerations have limited the choice of materials for the electrodes and cell walls to the various allotropic forms of carbon.\nTypically the Hall cell is a shallow vessel, with the floor forming the cathode, the side walls a rammed coke-pitch mixture, and the anode a block suspended in the bath at an anode-cathode separation of a few centimeters. The anode is typically formed from a pitch-calcined petroleum coke blend, prebaked to form a monolithic block of amorphous carbon. The cathode is typically formed from a pre-baked pitch-calcined anthracite or coke blend, with cast-in-place iron over steel bar electrical conductors in grooves in the bottom side of the cathode.\nDuring operation of the Hall cell, only about 25% of the electricity consumed is used for the actual reduction of alumina to aluminum, with approximately 40% of the current consumed by the voltage drop caused by the resistance of the bath. The anode-cathode spacing is usually about 4-5 cm., and attempts to lower this distance result in an electrical discharge from the cathode to the anode through aluminum droplets.\nThe molten aluminum is present as a pad in the cell, but is not a quiescent pool due to the factors of preferential wetting of the carbon cathode surface by the cryolite melt in relation to the molten aluminum, causing the aluminum to form droplets, and the erratic movements of the molten aluminum from the strong electromagnetic forces generated by the high current density.\nThe wetting of a solid surface in contact with two immiscible liquids is a function of the surface free energy of the three surfaces, in which the carbon cathode is a low energy surface and consequently is not readily wet by the liquid aluminum. The angle of a droplet of aluminum at the cryolite-aluminum-carbon junction is governed by the relationship ##EQU1## where .alpha..sub.12, .alpha..sub.13, and .alpha..sub.23 are the surface free energies at the aluminum carbon, cryolite-carbon, and cryolite-aluminum boundaries, respectively.\nIf the cathode were a high energy surface, such as would occur if it were a ceramic instead of carbon, it would have a higher contact angle and better wettability with the liquid aluminum. This in turn would tend to smooth out the surface of the liquid aluminum pool and lessen the possibility of interelectrode discharge allowing the anode-cathode distance to be lowered and the thermodynamic efficiency of the cell improved, by decreasing the voltage drop through the bath.\nTypically, amorphous carbon is a low energy surface, but also is quite durable, lasting for several years duration as a cathode, and relatively inexpensive. However, a cathode or a TiB.sub.2 stud as a component of the cathode which has better wettability and would permit closer anode-cathode spacing could improve the thermodynamic efficiency and be very cost-effective.\nSeveral workers in the field have developed refractory high free energy material cathodes. U.S. Pat. No. 2,915,442, Lewis, Dec. 1, 1959, claims a process for production of aluminum using a cathode consisting of the borides, carbides, and nitrides of Ti, Zr, V, Ta, Nb, and Hf. U.S. Pat. No. 3,028,324, Ransley, Apr. 3, 1962, claims a method of producing aluminum using a mixture of TiC and TiB.sub.2 as the cathode. U.S. Pat. No. 3,151,054, Lewis, Sept. 29, 1964, claims a Hall cell cathode conducting element consisting of one of the carbides and borides of Ti, Zr, Ta and Nb. U.S. Pat. No. 3,156,639, Kibby, Nov. 10, 1964, claims a cathode for a Hall cell with a cap of refractory hard metal and discloses TiB.sub.2 as the material of construction. U.S. Pat. No. 3,314,876, Ransley, Apr. 18, 1967, discloses the use of TiB.sub.2 for use in Hall cell electrodes. The raw materials must be of high purity particularly in regard to oxygen content, Col. 1, line 73-Col. 2, line 29; Col. 4, lines 39-50, Col. 8, lines 1-24. U.S. Pat. No. 3,400,061, Lewis, Sept. 3, 1968 discloses a cathode comprising a refractory hard metal and carbon, which may be formed in a one-step reaction during calcination. U.S. Pat. No. 4,071,420, Foster, Jan. 31, 1978, discloses a cell for the electrolysis of a metal component in a molten electrolyte using a cathode with refractory hard metal TiB.sub.2 tubular elements protruding into the electrolyte. The protruding elements enhance electrical conductivity and form a partial barrier to the mechanical agitation caused by magnetic effects."}
{"text": "In the course of processing substrates, for example semiconductor wafers or other kinds of substrates, photolithography is often performed to provide a mask for processing steps like etching or other operations which are to be applied only to part of the substrates. Typically, for photolithography a photoresist is applied to the substrate, the photoresist is exposed through a mask such that only defined regions of the photoresist are exposed to light, and the photoresist is subsequently developed. In case of positive photoresists, the regions exposed to light may be removed during the development, while for negative photoresists the regions not exposed to light are removed during the development of the photoresist.\nIn typical processing environments, substrates are transferred between the various processing steps like applying photoresist, exposing the substrate to light and developing the photoresist, by means of a mechanical gripper arm or other similar device which grips the substrate at a periphery thereof. However, this may lead to problems in some cases when the gripping device grips the periphery of the substrate where photoresist is applied, as the photoresist may for example adhere to the gripping device.\nIn some conventional solutions, the photoresist may be removed at the peripheral portion of the substrate prior to transferring the substrate using the gripping device. However, at least for some kinds of processing it may be desirable that the peripheral portion of the substrate remains coated with photoresist."}
{"text": "A common problem found in many trucks today is that their wheels and axles are incorrectly positioned. This incorrect position is too small for the human eye to see and the effects become noticeable only when the problem has become severe. The incorrect position is called wheel misalignment. Wheel misalignment leads to increased tire wear and fuel consumption and decreased control of the vehicle. Wheel misalignment can be due to components of a new truck setting in, bumps, potholes, suspension fatigue, or loading of the vehicle. It is therefore necessary to perform regular checks on the health of the vehicle.\nIf the truck owner decides to correct the problem of wheel misalignment, he has to take the vehicle to a service centre. The correction is done while the vehicle is at a standstill and without the extra weight of a cargo. The position of the wheel is then measured using alignment equipment. The position of the wheel is then altered by a technician till the values are within the recommended settings. Companies that perform wheel alignment commonly have a yearly check-up of their fleet. Midway the year, the truck wheels might already be misaligned, while the symptoms are too small to be noticed. This means that the truck might need to wait up to a year to be corrected.\nThis brings the need for onboard measurement equipment or dynamic alignment, to enable fleet operators to implement more effective maintenance programmes. The misalignment will be detected in an early phase, before causing too much damage to the tire. Dynamic alignment takes into consideration the loading pattern and dynamic behaviour of the vehicle whilst on the road, giving a more accurate measurement of the alignment angles."}
{"text": "1. Field of the Invention\nThe present invention relates generally to heat exchangers, and more particularly, to a heat exchanger for use in an automotive air conditioning system.\n2. Description of the Prior Art\nOne prior art embodiment of a heat exchanger as described in Japanese Patent Application Publication No. 2-154986 is essentially illustrated in FIGS. 1 and 2. As shown in the figures, a heat exchanger such as a condenser 100 includes a plurality of adjacent, essentially flat tubes 100 having oval cross-sections and open ends which allow refrigerant fluid to flow therethrough. A plurality of corrugated outer fin units 120 are fixedly disposed between adjacent tubes 110. Cylindrical header pipes 130 and 140 having closed ends are disposed perpendicular to flat tubes 110. Flat tubes 110 having open ends are fixedly and hermetically connected to the inside of header pipes 130 and 140, so as to communicate with the hollow interiors of header pipes 130 and 140.\nInlet pipe 131 has an open end which is fixedly and hermetically connected to the outside of an upper portion of header pipe 130. The other open end of inlet pipe 131 is linked to an outlet of an element (not shown) positioned upstream with respect to condenser 100. The element may be, for example, a compressor. Outlet pipe 141 also has an open end which is fixedly and hermetically connected to the outside of a lower portion of header pipe 140. The other open end of outlet pipe 141 is linked to an inlet of an element (not shown) positioned downstream with respect to condenser 100. This element could be, for example, a receiver.\nWith reference to FIG. 2, each flat tube 110 includes flat tube member 111 and a plurality of partition walls 112. Partition walls 112 are integrally formed along an inner surface of flat tube member 111. Partition walls 112 extend longitudinally along the length of flat tube member 111 so as to divide the interior hollow portion of flat tube member 111 into a plurality of rectangular parallelepiped hollow regions 113 and a pair of semicylindrical hollow regions 114 which are located at the lateral ends of flat tube member 111. Hollow regions 113 and 114 extend in parallel directions with respect to one another. As discussed below, the hollow regions extend transversely relative to a flow direction \"A\" of the air. The air flows along the exterior surface of the flat tube 110.\nDuring operation of a refrigerant circuit which includes condenser 100, the discharged refrigerant gas from a compressor is directed into the hollow interior of header pipe 130 via inlet pipe 131. The refrigerant gas directed into the hollow interior of header pipe 130 flows through the hollow interior of header pipe 130 toward its lower end. The refrigerant gas flowing through the hollow interior of header pipe 130 concurrently flows into each of the hollow regions 113 and 114 of each of flat tubes 110. The gas then longitudinally flows through each of hollow regions 113 and 114 of each of the flat tubes 110 from the right to the left sides (in FIG. 1). The refrigerant gas exchanges heat with air passing along corrugated fins 120 so as to be liquefied. The flow direction of the air passing along corrugated fins 120 is shown by large arrow \"A\" in FIG. 2. Accordingly, the air laterally passes along an exterior surface of flat tubes 110. Finally the refrigerant flows out from each of the hollow regions 113 and 114 of each of flat tubes 110. The liquefied refrigerant flowing out from each of hollow regions 113 and 114 of each flat tube 110 joins together at the hollow interior of header pipe 140, and flows through the hollow interior of header pipe 140 toward a lower end of header pipe 140. The liquefied refrigerant flowing through the hollow interior of header pipe 140 is conducted to the receiver via outlet pipe 141.\nIn this prior art embodiment, the integral formation of the partition walls 112 prevents expansion of flat tube members 111 caused by the pressure force of the refrigerant. Further, the area of the contact surface between the refrigerant and the flat tube 110 is increased so that a heat exchange efficiency of condenser 100 is improved.\nFIG. 3 essentially illustrates one of a plurality of identical flat tubes 210 which form a part of a condenser (not shown) as described in U.S. Pat. No. 4,998,580. As shown in the figure, flat tube 210 includes flat tube member 211 and corrugated inner fin 212. Inner fin 212 is fixedly disposed along the entire interior length of the hollow portion of flat tube member 211. Corrugated inner fin 212 includes a plurality of ridges 212a which longitudinally extend along the length of flat tube member 211. Adjacent ridges 212a are fixedly connected to upper and lower inner surfaces of flat tube member 211, respectively, so as to define a plurality of hollow regions 213. The hollow regions have a lateral cross section which is defined by the generally sine curve shape of fin 212, in the interior hollow portion of flat tube member 211. Hollow regions 213 are aligned with one another in a parallel relationship. Hollow regions 213 are oriented to extend transversely to the flow direction \"A\" of the air. The air passes along the exterior surface of the flat tube 210.\nIn this prior art embodiment, the attachment of the separately formed corrugated inner fin 212 to flat tube member 211 effectively prevents the expansion of flat tube member 211 caused by the pressure of the refrigerant, without unnecessarily increasing the thickness of flat tube member 211. Furthermore, the heat exchangeability of the condenser is improved because the refrigerant contacts a surface of corrugated inner fin 212 which conducts the heat to flat tube member 211.\nIn the prior art embodiments, the amount of heat exchange between the refrigerant and the air at the upstream side (relative to the air flow) the flat tube is large, while the amount of heat exchange between the refrigerant and the air at the downstream side (relative to air flow) region of the flat tube is small. This difference results because the refrigerant flows only longitudinally through each of the hollow portions of the flat tube; that is, the refrigerant does not disperse in the lateral direction when it flows through the interior hollow space of the flat tube. Therefore, the amount of heat exchange between the refrigerant flowing through the interior hollow space of the flat tube and the air laterally passing along the exterior surface of the flat tube is gradually decreased in the direction from the upstream side (relative to the air flow) region of the flat tube to the downstream side (relative to the air flow) region of the flat tube. Accordingly, the heat exchange efficiency of the condenser is not sufficiently improved by the divided construction."}
{"text": "1. Technical Field\nThe present disclosure relates to technique for extracting a target hair follicle from a donor site. More particularly, but not exclusively, the present disclosure relates to extracting the hair follicle from underneath an external surface of the skin.\n2. Discussion of Related Art\nHair restoration may be carried out by extracting hair follicular units from the patient's donor area and transplanting them in the patient's recipient area, which may be a thinning area. Generally, the donor area is at the back and sides of the scalp.\nConventionally, a strip of tissues may be removed from the donor area. The strip may be then dissected into individual follicular units. The units may then be transplanted into the recipient area. It has been observed that application of this technique results in a linear scar formed at the donor area.\nIn another conventional technique, a hair follicle is extracted by punching a hole around the hair follicle, and then extracting the hair follicle. The hole is punched through the external surface of the skin. It has been observed that application of this technique may result in dotted scars formed at the donor area. Further, in this technique, the hair follicle may be transected, as the punch may not be aligned with the alignment of the hair follicle.\nIn light of the foregoing discussion, there may be a need to reduce scarring in the donor site, and reduce transection or denuding of hair follicle during extraction.\nAfter the preparation of follicular unit grafts doctors make tiny holes in the patient's scalp at the recipient area where grafts are placed. The positioning and arrangement of follicular units depends upon aesthetic qualities of a hair transplant, case-by-case basis, depending on the patients' history of hair loss and likelihood of future hair loss.\nFrancisco Jimenez, MD, Ander Izeta, PhD, and Enrique Poblet, MD. “Morphometric Analysis of the Human Scalp Hair Follicle: Practical Implications for the Hair Transplant Surgeon and Hair Regeneration Studies” Dermatol Surg 4021; 37:58-64. This document is incorporated herein by reference in its entirety."}
{"text": "The present invention relates to progesterone receptor modulators.\nIntracellular receptors (IR) form a class of structurally related gene regulators known as “ligand dependent transcription factors”. The steroid receptor family is a subset of the IR family, including progesterone receptor (PR), estrogen receptor (ER), androgen receptor (AR), glucocorticoid receptor (GR), and mineralocorticoid receptor (MR). A compound that binds to an IR and mimics the action of the natural hormone is termed an agonist, whilst a compound that inhibits the effect of the hormone is an antagonist.\nThe natural hormone, or ligand, for the PR is the steroid progesterone, but synthetic compounds, such as medroxyprogesterone acetate or levonorgestrel, have been made which also serve as ligands. Once a ligand is present in the fluid surrounding a cell, it passes through the membrane via passive diffusion, and binds to the IR to create a receptor/ligand complex. This complex binds to specific gene promoters present in the cell's DNA. Once bound to the DNA the complex modulates the production of mRNA and protein encoded by that gene.\nPR agonists (natural and synthetic) are known to play an important role in the health of women. PR agonists are used in birth control formulations, typically in the presence of ER agonists, alternatively they may be used in conjunction with PR antagonists. ER agonists are used to treat the symptoms of menopause, but have been associated with a proliferative effect on the uterus that can lead to an increased risk of uterine cancers. Co-administration of a PR agonist reduces or ablates that risk.\nU.S. Pat. No. 6,407,101, which is hereby incorporated by reference, describes the preparation of cyclocarbamate derivatives, which are useful as progesterone receptor modulators. These cyclocarbamate derivatives, including, e.g., 5-(4,4-dimethyl-2-thioxo-1,4-dihydro-2H-benzoxazin-6-yl)-1-methyl-1H-2-cyano-pyrrole, are prepared by thionation of the corresponding benzoxazin-2-one (Scheme 1).\n\nWhat is needed in the art are alternate compounds that are effective as progesterone receptor modulators."}
{"text": "The sizes and shapes of automobiles and the types and locations of parts projecting therefrom, such as bumpers and the like, are constantly changing so that rotary brushes used in automatic equipment for washing vehicles are frequently and seriously damaged under circumstances which cannot be adequately anticipated. Thus, it is of great importance to provide a brush construction which can be quickly and easily, hence inexpensively, repaired to minimize loss of money and time, thereby increasing the efficiency in, and profit from, the car washing operation.\nIt is equally important that the rotary brushes be capable of quick and easy removal and replacement where such is necessitated by the wear and tear of the brush bristles resulting from ordinary use. Further, due to the size and shape of automobiles, and the various parts which project therefrom, selected areas of the brushes are subject to a more rapid rate of wear than are other brush areas. These selected areas, when worn, prevent adequate washing of the vehicle and thus it is often necessary to replace the complete brush assembly or, in the alternative, substantially disassemble and reassemble the complete brush assembly in order to replace a worn section.\nFurther, with the advent of rubber or other shock absorbing bumpers on the vehicles, it has been discovered that those areas on the brushes which contact the bumpers not only tend to have greatly increased wear, but also have a tendency to \"hang-up\" or become locked to the car bumper. This can result in damage to the vehicle, brush or both.\nThe increased wear caused by vehicle projections, such as bumpers, obviously requires substantially increased maintenance and replacement of the brushes. Likewise, the potential damage which worn brushes can cause to the vehicles has increased the maintenance required and has thus decreased the profit of such car-washing operations.\nMany attempts have been made to produce brushes which would successfully overcome the above-mentioned disadvantages, and one such brush construction is disclosed in my prior U.S. Pat. No. 3,439,373. However, in a continuing effort to improve upon, and reduce the cost of, the procedure of replacing worn or damaged brushes or brush bristles, while at the same time increase brush wear, I found that much time could be saved by constructing the brush from a plurality of removable arcuate brush sectors which could also have the bristle density thereof varied as desired to result in optimum washing and wear characteristics relative to the configuration of the vehicles.\nAccordingly, a primary object of this invention is the provision of an improved rotary brush structure which provides a rotary drumlike hub or frame having a plurality of removable arcuate brush sections mounted thereon. The brush sections preferably comprise semi-cylindrical shells which are mounted in opposed pairs on the hub. Several pairs may be disposed axially adjacent one another along the length of the hub. This brush structure thus enables individual worn sections to be easily removed and replaced without requiring disassembly of the entire brush construction. At the same time, this structure permits various brush sections to be axially interchanged or their orientation reversed to provide for more uniform wear on all of the brush sections.\nA further object of the present invention is to provide an improved brush construction, as aforesaid, which enables the arcuate brush sections to be attached to or removed from the hub in a simple manner while at the same time the arcuate brush sections can be individually economically constructed.\nStill a further object of the present invention is to provide an improved brush construction, as aforesaid, which is extremely durable but of rather light weight so as to facilitate the mounting and driving thereof to thus minimize wear on the bearings and other associated driving parts.\nOther objects and purposes of the invention will be apparent to persons familiar with rotary brush constructions upon reading the following description and examining the accompanying drawings."}
{"text": "This invention relates to localized thermal processing of glazed ceramics and other brittle and low thermal conductivity materials. More particularly, this invention relates to coating and filler materials for use in such processing, the coating materials having coefficients of thermal expansion of less than about 8xc3x9710xe2x88x926/xc2x0 C. and glass transition temperatures of less than about 400xc2x0 C., and the filler materials having coefficients of thermal expansion of about 0.5xc3x9710xe2x88x926/xc2x0 C.\nVitreous coating materials, commonly referred to as glazes or enamels, are used on many ceramics and other brittle and low thermal conductivity materials for aesthetic and functional purposes. Ceramic whitewares are coated with glazes for decorative and sanitary purposes. Conductive glazes are used on ceramic electrical insulators to remove charges. Glazes are also used in ceramic electronic packages for hermetic sealing and bonding. Enameled cast iron is coated with porcelain enamels for decorative and functional purposes.\nGlazed ceramic manufacturing processes often yield a certain percentage of product that have surface defects. These defects can detract from the appearance of the article and create functional problems that render the article unusable. Consequently, methods have been developed for repairing certain types of surface defects in glazed ceramics. One such repair method involves enlarging the hole by drilling, filling the hole with a glaze material similar to the original, and re-firing the entire part at maximum temperatures near 1200xc2x0 C. for many hours. Unfortunately, new flaws may initiate at secondary sites during re-firing, thus, frustrating the repair efforts. Moreover, this method is expensive, energy intensive, and time consuming.\nAnother repair method considered for use in repairing glazed ceramics is laser welding. Laser welding methods typically employ a carbon dioxide (CO2) laser operating at a wavelength of 10.6 xcexcm. Such lasers are used because they are stable and capable of delivering high, readily absorbed output power such that temperatures of 1800xc2x0 C.-1900xc2x0 C. over a few square millimeter area can be rapidly achieved to fuse refractory materials within a few seconds. The rapid heating can also produce undesirable cracking within and near the irradiated region due to the brittle and low thermal conductive nature of ceramics which permit high thermal gradients and high stresses to develop.\nLaser induced thermal stress problems have been addressed in U.S. Pat. No. 5,427,825 which discloses a method for localized glazing of ceramics that maintains a non-destructive level of thermal stresses within the irradiated region by appropriate control of the laser\"\"s energy. This is accomplished by applying a fusing radiant energy to a fusion zone on the surface of the glazed ceramic article and annealing the fusion zone by applying annealing radiant energy to an annealing zone on the surface of the article which encompasses the fusion zone. The annealing radiant energy limits the thermal stress to a level during cooling which is less than the fracturing stress of the ceramic material, so that cooling can be accomplished without inducing cracks in the surface.\nThe glazes coating the ceramics described in U.S. Pat. No. 5,427,825 typically have melting temperatures near 1200xc2x0 C., and coefficients of thermal expansion of 6-7xc3x9710xe2x88x926/xc2x0 C. The engineering of these two parameters has been found to be critical in laser fusion performance, particularly in laser sealing of vitreous coatings over brittle substrates. The high melting temperature requires the absorption of considerable energy to fuse the coating to the substrate. Consequently, the non-destructively high thermal gradients and the thermal expansions in the range noted, can still increase thermally-induced strains beyond critical limits, thus resulting in permanent damage to the coating and the substrate.\nAccordingly, coating materials are needed which have coefficients of thermal expansion and glass transition temperatures sufficiently lowered to a point where the strain induced by a rapid change in temperature during localized thermal processing, remains below the elastic strain limit. There is also a need for filler materials that can be used for repairing deep defects via localized thermal processing.\nThe present invention relates to coating and filler materials for localized thermal processing of glazed ceramics and other brittle and low thermal conductivity materials. The coating materials comprise oxide compositions that exhibit coefficients of thermal expansion which are less than about 8xc3x9710xe2x88x926/xc2x0 C. and glass transition temperatures which are less than about 400xc2x0 C. The filler materials comprise particulate oxide materials which do not substantially react during localized thermal processing of glazed ceramics and other brittle and low thermal conductivity materials.\nThe coating and filler materials are useable together as a composite material for repairing cavities having depths greater than about 2 mm."}
{"text": "The present technology relates to an image processing apparatus, an image processing method and electronic equipment. More particularly, the present technology relates to an image processing apparatus capable of preventing the image quality from deteriorating and relates to an image processing method provided for the apparatus as well as electronic equipment employing the apparatus.\nIn the past, a solid-state imaging device represented by a CMOS (Complementary Metal Oxide Semiconductor) image sensor has always been required to have a smaller pixel size and a larger pixel count representing the number of pixels which can be provided in the same image area.\nThe configuration of an ordinary solid-state imaging device is explained by referring to FIGS. 1A to 1C as follows.\nFIG. 1A shows a typical color filter 11 for a solid-state imaging device whereas FIG. 1B shows a cross section of a front-irradiation CMOS image sensor 21. FIG. 1C shows a cross section of a rear-irradiation CMOS image sensor 31.\nThe color filter 11 shown in FIG. 1A includes red, blue and green color filters which are laid out to in the so-called Bayer array. In the Bayer array, a number of pixels each having a cell positioned to face one of the color pixels are laid out in the horizontal and vertical directions. In the Bayer array, B (blue) and Gb (green) color filters are laid out alternately in the horizontal direction along a row so that each of the color filters faces the cell of one of pixels on the same row. By the same token, in the Bayer array, R (red) and Gr (green) color filters are laid out alternately in the horizontal direction along another row so that each of the color filters faces the cell of one of pixels on the same other row. That is to say, the pixels whose cells face the B (blue) and Gb (green) color filters are laid out in the horizontal direction in the Bayer array every other row whereas the pixels whose cells face the R (red) and Gr (green) color filters are laid out in the horizontal direction in the Bayer array also every other row. In addition, the pixels whose cells face the B (blue) and R (red) color filters are not laid out on the same column in the vertical direction.\nIt is to be noted that, in the following description, a pixel whose cell faces an R (red) color filter is referred to as an R pixel whereas a pixel whose cell faces a Gb (green) color filter is referred to as a Gb pixel. By the same token, a pixel whose cell faces a B (blue) color filter is referred to as a B pixel whereas a pixel whose cell faces a Gr (green) color filter is referred to as a Gr pixel.\nThe front-irradiation CMOS image sensor 21 shown in FIG. 1B is configured to have a silicon substrate 22 including a photodiode. In addition, the front-irradiation CMOS image sensor 21 also includes an FD (Floating Diffusion) 23 and a reflection prevention film 24. The reflection prevention film 24 is created on the silicon substrate 22 and the FD 23. A wiring layer 26 having wires 25 is created on the reflection prevention film 24. A flattening film 27 is created on the wiring layer 26 whereas a color filter 28 is created on the flattening film 27. Then, an on-chip lens 29 is provided on the color filter 28.\nThe rear-irradiation CMOS image sensor 31 shown in FIG. 1C is configured to have a silicon substrate 32 including a photodiode. In addition, the rear-irradiation CMOS image sensor 31 also includes a reflection prevention film 33 created on the silicon substrate 32. A light shielding film 34 for preventing crosstalk is created on the reflection prevention film 33. A flattening film 35 is created on the light shielding film 34 whereas a color filter 36 is created on the flattening film 35. Then, an on-chip lens 37 is provided on the color filter 36. It is to be noted that, in the case of the rear-irradiation CMOS image sensor 31, the on-chip filter 37 for receiving incident light is provided on the rear side and a wiring layer not shown in FIG. 1C is provided on the front side.\nIn the rear-irradiation CMOS image sensor 31, the wiring layer is not provided on the light-incidence side. Thus, incident light is not lost due to the wiring layer so that the amount of incident light arriving at the silicon substrate 32 can be increased to a quantity greater than that of the front-irradiation CMOS image sensor 21. As a result, by making use of the rear-irradiation CMOS image sensor 31, it is possible to obtain high-sensitivity, low-noise and high-quality image. Such rear-irradiation CMOS image sensors 31 are mass-produced and employed in electronic equipment such as a cam coder and a digital still camera.\nBy the way, since the number of pixels employed in a solid-state imaging device is increased, the absolute quantity of the energy of light incident to a pixel undesirably decreases and crosstalk inevitably occurs. In this case, the crosstalk is a phenomenon in which light leaks out to an adjacent pixel existing among pixels while the light is propagating through the device employing the pixels. In addition, the number of electrons obtained as a result of opto-electrical conversion taking place in the neighborhood of a pixel boundary rises, unavoidably increasing crosstalk as well. In this case, the crosstalk is a phenomenon in which electrons leak out to an adjacent pixel. As a result, these kinds of crosstalk increase. The generation of these kinds of crosstalk is a cause deteriorating the spectroscopic characteristic of the rear-irradiation CMOS image sensor 31.\nNext, the spectroscopic characteristics of the rear-irradiation CMOS image sensor 31 are explained by referring to FIGS. 2A to 2C.\nFIG. 2A shows spectroscopic characteristics found from signals output by a rectangular rear-irradiation CMOS image sensor 31 having a pixel size of 1.12 micrometers. In FIG. 2A, the horizontal axis represents the wavelength expressed in terms of nm whereas the vertical axis represents an output signal (arb. unit) which is the magnitude of a signal output by the rectangular rear-irradiation CMOS image sensor 31.\nAs shown in FIG. 2A, the pixel size has a value of an order not much different from the wavelength so that the color separation becomes poor.\nFIG. 2B shows spectroscopic characteristics for a configuration including an on-chip filter. In FIG. 2B, the vertical axis represents the transmission (arb. unit) whereas the horizontal axis represents the wavelength lambda expressed in terms of micrometers.\nThe spectroscopic characteristics shown in FIG. 2B to serve as the spectroscopic characteristics for a configuration including an on-chip filter are compared with the spectroscopic characteristics shown in FIG. 2A as spectroscopic characteristics found from signals output by a rectangular rear-irradiation CMOS image sensor 31 in order to clarify the following. The deterioration of the color separation does not depend on the characteristic of the on-chip color filter, but depends on leaks of light or electrons (that is, signal electric charges) inside the rear-irradiation CMOS image sensor 31.\nFIG. 2C shows spectroscopic characteristics found from signals output by a rectangular rear-irradiation CMOS image sensor 31 having a pixel size of 1.12 micrometers for a case in which the incidence angle of light incident to the light receiving surface of the rear-irradiation CMOS image sensor 31 is set at ten degrees. On the other hand, FIG. 2A described above shows spectroscopic characteristics found from signals output by a rectangular rear-irradiation CMOS image sensor 31 having a pixel size of 1.12 micrometers for a case in which the incidence angle of light incident to the light receiving surface of the rear-irradiation CMOS image sensor 31 is set at zero degree. By comparing the spectroscopic characteristics shown in FIG. 2A with those shown in FIG. 2C, it becomes obvious that an increase of the incidence angle of light incident to the light receiving surface of the rear-irradiation CMOS image sensor 31 emphasizes crosstalk.\nIn addition, for example, applicants for a patent of the present technology have also proposed an imaging apparatus capable of reducing effects on the image quality by carrying out pixel color mixing correction processing in accordance with a correction parameter (refer to, for example, Japanese Patent Laid-Open No. 2009-124282)."}
{"text": "The invention relates to a circuit board for connecting a deformation sensor, which is arranged on a radially outer side of a rolling bearing outer ring, to a signal-processing circuit. The invention also relates to a wheel bearing for a vehicle with the circuit board.\nFrom DE 101 64 929 B4 it is known to arrange deformation sensors circumferentially on a radially outer side of an outer ring of a wheel bearing, wherein these sensors detect deformation forces acting on the outer ring and indicate these forces in the form of electrically detectable signals. Based on these signals, the deformation forces on the outer ring can then be detected in all spatial directions, in order to use them, for example, for stabilizing the vehicle."}
{"text": "In the manufacture of semiconductor devices, a semiconductor wafer is divided by scribe lines into a plurality of dies. The dies comprise integrated circuits of identical type and, after appropriate testing to verify operability and reparability of the individual dies, they are sawed from the wafer and are thereafter referred to as chips. Each chip is appropriately packaged to result in the completed semiconductor device.\nEach die has a plurality of contact areas, which are typically referred to as pads. In the completed semiconductor device, such pads are connected to leads which are accessible via the device packaging. During electrical testing of the wafer, individual contact areas are contacted by respective testing probes of a probe card. Such probes are typically needle-like elements mounted on a probe card structure, which are aligned with the contact areas before actual contact is made. The needle-like devices penetrate slightly into the contact areas to make electrical contact. There are probe types other than needle-like devices which may be used with various types of pad structures, but the present invention is not dependent on the particular type of contact area or probe which is utilized.\nFor functional testing, dedicated automated test equipment (henceforth referred to as “IC test system”) is connected to the probe card and voltages are applied through the probes to various contact areas of the wafer. Basic chip operation usually requires a minimum of three types of pads: (1) pads through which a global supply voltage VS is applied, (2) pads that are used to control the command logic of the integrated circuit, and (3) pads that are used to supply data bits to and receive data bits from the chip. Those return signals from the wafer also pass through designated probes and are fed back to the electrical tester. The location of test failures is stored in a fail memory so that repairs can be kept track of, and the testing sequence is controlled automatically.\nA probe card may be designed so that its probes cover multiple dies in order to maximize the number of dies can be tested at the same time. The degree to which multiple dies can be tested simultaneously is referred to as parallelism. Each probe is ultimately connected to the electrical tester, and corresponds to a channel of the testing system. Since the factors limiting parallelism are the capacity of the fail memory and the number of available testing channels per device under test, it is desirable to maximize the number of channels which are available for testing.\nTypically, less than all of the contact areas on the dies are utilized for testing, and it may be a requirement of the test that one or more of the contact areas which is not tested be held at a predetermined voltage value during the test. In accordance with a prior mode of accomplishing this, the predetermined voltage value was obtained through the probe card, and since any probe(s) used for obtaining such value could not be used for testing, a limitation on the number of testing channels resulted."}
{"text": "Electronic data interchange (EDI) is one of the ways businesses use for exchanging computer-to-computer business information based on approved formatting standards and schemas. For example, millions of companies around the world transmit data associated with business transactions (e.g., purchase orders, shipping/air bills, invoices, or the like) using EDI to conduct commerce.\nIn a typical EDI transaction model, a large business entity or an EDI integration broker trades with numerous partners and has the technical capability to handle numerous EDI transaction data in various EDI formats and schemas. These large business entities, also known as “hubs,” transact with one or more suppliers, also known as “spokes”. Each of the spokes typically is a relatively small business entity that is only capable of dealing with one hub.\nFrequently, EDI documents are exchanged in batches between the hub and the spokes, and each batch may include multiple types or a collection of business documents. For example, the batch of EDI documents from the hub to the spokes may include: purchase order, payment advice, return material notifications etc.\nTrading parties also may modify the EDI formats, and consequently the schemas associated with the EDI documents, during their course of conducting business to serve their business and integration needs. After such modification, business entities frequently use an EDI transformation system to transform EDI data or EDI messages formatted from the previous schema to the modified schema.\nIn transforming EDI transaction data formatted according to previous or old EDI schemas, EDI transformation systems (e.g., systems using XSLT) typically perform a transformation of all EDI transaction data from the previous EDI schemas to the new EDI formats. As such, such systems load the EDI message, which includes a batched transaction data with a size of 30 to 40 MB, in a memory storage area for the transformation to take place. This transformation thus unnecessarily poses a limit on the size of documents that can be loaded without compromising performance significantly. In addition, for the spokes that lack advanced processing equipment to handle such transformation efficiently, existing transformation method places heavy burden on the transformation performance and resources of the spokes."}
{"text": "1. Field of the Invention\nThe present invention relates to an internal voltage dropping circuit for a semiconductor device, and more particularly, it relates to an internal voltage dropping circuit by which current consumption during a standby can be reduced.\n2. Description of the Prior Art\nIn recent years, semiconductor integrated circuits have been increasingly miniaturized. Especially, that is most noticeable in dynamic RAMs. As transistors are miniaturized more and more, supply voltage must be reduced for various reasons, such as lifetime shortening caused by hot electrons. For instance, in transistors of 0.6 .mu.m gate length, supply voltage applied from the outside thereto is 5 V but must be reduced to 4 V or below (e.g., 3.3 V) by use of an internal voltage dropping circuit.\nFIG. 8 shows an embodiment of a conventional internal voltage dropping circuit for an integrated circuit semiconductor load device. This circuit employs a current mirror type differential amplifier. The current mirror type differential amplifier 22 consists of two P-type MOSFETs 42 and 43 and two N-type MOSFETs 44 and 45, which are symmetrically arranged on two current paths, and one of the N-type MOSFETs, 44, has its gate connected to a reference potential terminal 24 while the other MOSFET 45 has its gate connected to an output terminal 25. A power source is connected via an output transistor 26 which is a P-type MOSFET to the output terminal 25, and a potential at a junction of the transistors 43 and 45 on the output path in the differential amplifier 22 is applied to a gate of the output transistor 26. There are two paths provided between the differential amplifier 22 and a lower potential power source (the ground in the embodiment in FIG. 8), and a transistor 27 which permits merely a small amount of current to flow is connected to one of the paths i.e. including transistor 44 while a transistor 28 which permits a relatively large amount of current to flow is connected to the other path i.e. including transistor 45.\nThe internal voltage dropping circuit works as follows. When a potential at the output terminal 25 is lower than that at the reference potential terminal 24, the differential amplifier 22 causes a gate potential at the output transistor 26 to drop, and the output transistor 26 turns on to supply current to the output terminal 25 until the output terminal 25 reaches the same potential as the reference potential terminal 24. The transistors 27 and 28 are for saving the current flowing in the differential amplifier 22, and when an active/standby switch signal 41 received from a control device (not shown) becomes low when a load, such as a CPU and the like, connected to the output terminal 25 is on standby, the transistor 28 turns off, and the transistor 27 permits merely a small current to flow. On the other hand, when the active/standby switch signal 41 becomes high when the load is activated, the transistor 28 turns on to permit a sufficient current to flow in the differential amplifier 22, and the operation speed of the differential amplifier 22 is enhanced.\nAs has been described, in the conventional internal voltage dropping circuit, even if the semiconductor device 28 is on standby, the transistor 27 permits a small current to always flow by virtue of the gate thereof being coupled to the same power source as transistors 26, 42 and 43, and this wastes electric power.\nA power supply circuit having a power-down mode for reducing a current supply to a SRAM as a semiconductor device is disclosed in 1987 IEEE International Solid-State Circuit Conference Digest of Technical Report pp 252-253, \"A 256 K SRAM with On-Chip Power Supply Conversion,\"."}
{"text": "1. Field of the Invention\nThe present invention relates to the manufacture of reinforced concrete. In particular the present invention relates to reinforced concrete wherein the rebar within the reinforced concrete is secured with a plastic tie fastener or a substantially similar device. Furthermore, the present invention relates to a method for manufacturing reinforced concrete using plastic tie fasteners or a substantially similar device. The present invention also relates to devices for securing rebar in reinforced concrete.\n2. Description of the Related Art\nThe process of manufacturing steel reinforced concrete begins by forming a two or three dimensional grid of steel rebar around which concrete is poured. The individual lengths of rebar 10 are secured to one another to form this grid. Currently the industry standard is that the lengths of rebar are secured to one another at the grid intersections either by wrapping metal (usually steel) wire 12 around all lengths of rebar which form the grid at the intersections or by securing the grid intersections with spring steel wire clips (see FIGS. 1a, 1b and 2).\nIf metal wires 12 are used to prepare the grid, the free ends of the metal wire are twisted to tension the tie, securing the crossed bars in place. The free ends of the tie wire are then cut off or bent away from the concrete surface above the twist.\nIf the spring steel wire clips are used, the intersections are held together by the spring force of the spring steel wire clips.\nIf the concrete product is to be placed in a hostile environment (such as salt water, basic or acidic soil, etc.), the additional step of coating the metal tie wires or spring steel wire clips with some sort of protective coating or of using stainless steel metal tie wires must be taken.\nThere are a few drawbacks of the present uses of steel wire 12 and/or steel spring clips. The first drawback is the added time required by a construction worker to twist the free ends of the metal wires to tension the ties. Furthermore, the twisting of the wire fatigues the metal wire 12 at the twist site. This could cause the wire to break and allow the rebar 10 to become unsecured either prior to the pouring of concrete or even worse after the concrete has been poured. If the rebar becomes unsecured after the concrete has been poured the concrete may not have the reinforced strength that is required.\nAnother drawback of using metal wire ties is that the act of the construction worker attaching a machine to twist the ends of the wire tie, or the act of twisting the wire tie by hand is a time consuming task. Furthermore, the act of twisting wires requires an amount of area next to the rebar to twist the wire. Many times there are so many rebars 10 used to reinforce the concrete that there is little room between adjacent rebars to twist the steel wire 12.\nAnother drawback, which was alluded to above, is that the steel wires 12 or clips will rust if the concrete product is placed in a hostile environment. The rusting may weaken the overall structure and or produce unsightly stains on the structure. Rusting of the steel wire or clips cannot be completely avoided, but can be inhibited to some degree by coating the wires or clips with a protective coating. The extra coating step is costly and not always successful due to incomplete coating of the wires, having a twisted end which is difficult to completely coat, and if the end is removed or breaks off during the process an exposed end is left.\nThe present invention relates to the manufacture of reinforced concrete structures wherein plastic tie fasteners are used to secure the rebar structure within the concrete structure. Furthermore, the present invention relates to a method of securing rebar by using plastic tie fasteners. The plastic tie fasteners are looped about the rebar bars and tightened using a unidirectional force. The plastic tie fasteners can be looped and tightened via hand or an applicator or tensioning apparatus. The plastic tie fasteners offer the unexpected advantages of strength, flexibility, and ease of use to the reinforced concrete industry. Furthermore, the plastic tie fasteners do not rust or require rust retardant coatings. Still furthermore, the present invention prescribes a rebar securing apparatus comprising a plastic strip with a means for securing an end of said plastic strip in a looped configuration. The securing apparatus\"\" means for securing can be a ratcheting or slip resistant receptacle. Furthermore, the means for securing can be designed to disallow the loop from becoming larger or can be designed to allow the loop to become larger without damaging the securing apparatus.\nThe present invention provides a multitude of advantages which include, but are not limited to being incapable of rusting and staining a concrete structure, decreasing the amount of time required to secure rebar bars, an ease of use due to a simple technique for tightening the loop of the plastic tie fastener about the rebar, requiring only one hand of a construction worker to attach the plastic tie fastener to the rebar, and decreasing construction costs."}
{"text": "The present invention relates generally to data processing. More particularly, the present invention relates to processing image data in mosaic form.\nA variety of devices such as mobile phones are now equipped with image sensors such as charge-coupled-devices (CCD). These devices comprise an array of elements, each of which is filtered so as to sense only one color of light, typically red, green or blue. For each element the sensor provides a pixel value representing the intensity of the color of light detected by that element. These pixel values are referred to as “raw data” or “mosaic data.” Because images represented by mosaic data (referred to herein as “mosaic images”) are not pleasing to the eye, they are “demosaiced” (that is, converted into RGB images) before display. Each pixel in an RGB image is represented by three pixel values, each representing the intensity of one of the respective colors.\nIn conventional devices, demosaicing is performed before any further image processing, such as color processing, compression and decompression, and image enhancement. One disadvantage of this approach is that, because the volume of RGB data is three times that of mosaic data, more resources are consumed when processing RGB data."}
{"text": "This invention relates to an information communicating apparatus connected to a telephone subscriber line, such as a modem device, a telephone set, a facsimile machine, etc., and, more particularly, to a line interface circuit.\nIn the field of communication, in order to provide for both protection of network facilities having high public use and protection of terminals, a high insulation quality is required at the boundary between a network and a terminal, and, for that reason, compact transformers for communication which make provision for a high insulation quality have been used previously. However, a problem has occurred in that, in proportion to the popularization and evolution of personal terminals, it is required for transformers which are further miniaturized and reduced in weight to be used for a portable terminal; however, this requirement can not be satisfied sufficiently only with improvement in the material and arrangement of parts in the transformers. In light of this, the possible use of an isolator has been examined.\nAlso, in the application of this technology to the field of measurement, medical treatment or the like, it may be necessary to isolate a signal detecting part and a signal handling part, e.g. a sensor and a signal processing circuit. It is known that an isolator can be used as an isolation and separation means in this case.\nFor the purpose of making the isolator miniature, reliable and low-cost, a capacitive isolator has been developed. A ceramic capacitor for power application and surge protection, which is made on the basis of a high voltage capacitor technique providing an individual part constituting an-isolation barrier, is known; and, a circuit block for signal transmission using this capacitor is referred to as a capacitive isolation amplifier or a capacitive isolation coupler. This device has been used Since the 1970s.\nAlso, what is disclosed in Japanese Patent Application Laid-open N6-.62-260408 (1987) achieves isolation and separation of an analogue signal transmission in which a pulse signal obtained by the modulation of the analogue signal is transmitted to the secondary side, which is insulated and separated by a capacitor; and, after decoding, the analogue signal is regenerated.\nFurther, Japanese Patent Application Laid-open No. 11;136293 (1999) entitled xe2x80x9cIsolator Circuit and Monolithic Isolatorxe2x80x9d and Japanese Patent Application Laid-Open No. 11-317445/1999 entitled xe2x80x9cIsolator and Modem Device using itxe2x80x9d disclose a technique by which an isolator is formed on an LSI so that it has an isolation and separation function. This technique also provides means for achieving a system LSI including an insulation and separation function for the modem, etc.\nA system LSI for a modem including the insulation and separation function according to the above-mentioned technology dispenses with an insulation transformer and a photo coupler, which were needed for insulation previously, whereby it is possible to reduce the number of parts, the cost and the size of the mounting area. It also enables more suitable system design by improving the degree of freedom for chip division without encountering a restriction of the insulating location.\nAlso, by using electric power supplied from the system switch to operate the line side of two regions (line side: subscriber line side, host side: microprocessor side) on the LSI insulated by an isolator, it is possible to reduce the power consumption of the host side by a large amount. Therefore, it is possible to extend the life of a battery driving the host side in a portable terminal.\nHowever, in the system LSI for a modem including the insulation and separation function disclosed in the above-mentioned Japanese Patent Application Laid-open No. 11-136293 (1999) and Japanese Patent Application Laid-Open No. 11-317445 (1999), further consideration needs to be given to supplying power to a circuit used at the time of the reception of the caller identification information (Caller ID).\nIn the above-mentioned prior art, the supply of power from the system switch uses a loop current which determines whether a telephone set is in an on-hook condition or an off-hook condition. In general, it is required that a loop current of 20 to 120 mA is caused to flow to inform the system switch of the fact that the telephone set is in an off-hook condition, depending upon the standards defined by the government. Also, in order to inform the system switch of the fact that the telephone set is in the on-hook condition, the loop current must be maintained below a predetermined constant value (for example, 3 mA). Further, at the time of the reception of the Caller ID indication, the on-hook condition must be established in many countries (for example, in Japan), and a circuit for the reception of the Caller ID indication must operate with the consumption of electric power below a predetermined constant value (for example, 3 mA).\nIn order to solve the above-mentioned problems, the present invention utilizes the following means:\n(1) An isolator interface has a parallel/serial conversion circuit for converting the output of a digital analogue converter having plural bits to a serial signal, or a multiplexer for transmitting in a time divisional manner the output of the analogue digital converter and the status information of a line interface circuit. At the time of a normal operation, the output of the analogue digital converter is input to the isolator via the isolator interface, and at the time of the reception of a call signal, or at the time of the reception of a caller identification signal, the output of the analogue digital converter is input directly to the isolator.\n(2) The isolator interface has a parallel/serial conversion circuit for converting the output of the analogue digital converter having plural bits to a serial signal, or a multiplexer for transmitting in a time divisional manner the output of the analogue digital converter and the status information of the line interface circuit, and a clock generating circuit for producing a clock signal of a first frequency and a clock signal of a second frequency lower than the first frequency. At the time of normal operation, the analogue digital converter and the isolator interface are driven by the clock signal of the first frequency output from the clock generating circuit, and, at the time of the reception of the call signal, or at the time of the reception of the caller identification information, the analogue digital converter and the isolator interface are driven by the clock signal of the second frequency output from the clock generating circuit."}
{"text": "The electronic changes which occur in chemical or biological molecules when interacting with other molecules provide a sensitive and convenient method of studying molecule/molecule association. For example, technology capable of recording the electronic changes which occur in specific DNA/RNA sequences has been of great importance in numerous applications, such as in medical research and clinical diagnosis.\nIn general, DNA/RNA array sensors are based on single-stranded oligonucleotide (ODN) probes (of known sequence) immobilised on a physical substrate.\nAddition of tagged ODNs in the sample may result in immobilisation of tagged ODNs via Watson-Crick base pairing (or hybridization) if the sequences are complementary. The resulting double stranded ODN is immobile, being anchored by the probe, as well as tagged by prior sample tagging.\nThe immobilisation of tagged ODNs provides the basis of sample detection and the resulting pattern is detected by (for example) fluorescence or chemiluminescences from the tag. The main disadvantages of such detection methods are cost, complexity and inefficiency of tagging reactions.\nA number of alternative ‘label-free’ methodologies have also been developed, such as acoustic (Okahata et al., 1992; Ma, et al., 2002), optical (Piunno et al., 1994; Isola et al., 1998) and electrochemical (XU et al., 2000; Wang et al., 2001, 2003; Hashimoto et al., 1994; Napier et al., 1997; Pividori et al., 2000) methods which seek to detect the hybridization event directly. Electrochemical approaches using metal complexes (Takenaka et al., 2000), organic redox indicators (Millan et al., 1993) or nano particles (Wang et al., 2001, 2003) have also been investigated for their suitability to report hybridization.\nConducting polymers are attractive substrates for nucleotide sensors as they can act as an electronic transducer for charged species binding to their surface, and several methods for the immobilisation of ODN probes onto conducting polymers have been reported. In addition, the rich chemistry associated with fabrication of conducting polymer sensors and electrodes allows the detection of a broad range of chemical or biological materials according to the selection of probe and the capacity of the conducting polymer to bind the probe (Janata & Josowicz, 2003; Saxena & Malhotra, 2003).\nInitial approaches to construct conducting polymer ODN sensors included direct adsorption of the ODN probes onto oxidized polypyrrole (PPy) films by electrostatic attractions (Palecek, 2002; Minehan et al., 1994) or incorporation of the ODN probe into the polymer film as a macro-counterion (Wang et al., 1999). However, these methods constrain the orientation of the ODN probe resulting in high steric and kinetic barriers to hybridization, and the possibility of oxidative damage to the ODN probes. The result is that such methods have poor sensory properties.\nTo overcome steric constraints, several groups have electro copolymerised N-position substituted pyrroles with pyrrole to functionalise the polymer for ODN probe linkage to the substrate (Wang et al., 2000; Livache et al., 1995; Lassalle et al., 2001). However, the resulting poly(N-substituted pyrrole) possessed lower conductivity, possibly due to a lack of ring planarity in the resulting polymer.\nOf most interest has been a report that a 3-substituted pyrroles gave rise to a polymer film with higher conductivity (Delabouglise et al., 1989). A probe may be chemically coupled to a polymer by an amino-modified ODN.\nU.S. Pat. No. 6,096,825 discloses such a conducting polymer formed from a 3-substituted polypyrrole polymer. The conducting polymer includes a linker unit located in the 3-position of pyrrole monomers. The linker is bound to the pyrrole, at one end, and to the probe at the other end.\nThe linker connecting the probe to the polymer represents a key feature in electrochemical transduction between the probe and the polymer. Typically short saturated organic chains of 2- to 4-carbon atoms have been used to connect the probe to the polymer (see also H. Peng et al Label-free electrochemical DNA sensor based on functionalised conducting polymer Biosensor and Bioelectronics 20 (2005) 1821-1828).\nThere is a continuing need to synthesize new conducting polymers, and in particular linkers which enhance or improve electronic transduction between the probe and the polymer backbone. In addition, there is a continuing need to provide conducting polymers which have negligible effect on the chemical or biological sensory properties of probes they are attached to.\nThroughout this specification reference is made to “unsaturated” or “conjugated” organic chains. The words “unsaturated” and “conjugate” refer to double or triple bonds between two or more atoms. The words are used interchangeably throughout the specification.\nThroughout the specification reference is made to “probes”. The use of the word “probe” in this specification is intended to be interpreted inclusively and including any chemical or biological species which may be tethered via a linker to a polymer backbone. A probe may include a single or double stranded ODN, a virus, a protein, a cell, a peptide nucleic acid (PNA), a polysaccharide, a drug, or an RNA, for example.\nThroughout the specification reference is made to the phrase “sensory properties” in relation to a probe interaction with other chemical or biological molecules or species. Sensory properties should be interpreted as the ability of the probe to hybridize and/or otherwise interact through non-covalent bonding or the like with another chemical or biological molecules or species and to be efficiently transduced into a useful signal. Ideally a biological probe, for example, which is bound to a conducting polymer as far as possible, retains its natural sensory properties."}
{"text": "The present invention relates generally to integrated circuit chip package technology, and more particularly to a unique package design that facilitates the efficient placement of the leads therein.\nIntegrated circuit dies are conventionally enclosed in plastic packages that provide protection from hostile environments and enable electrical interconnection between the integrated circuit die and an underlying substrate such as a printed circuit board (“PCB”).\nSemiconductor packages employing micro lead frames (“MLF”) typically include a metal lead frame (often made from copper (Cu)), an integrated circuit die, bonding material to attach the integrated circuit die to the lead frame, bond wires which electrically connect pads on the integrated circuit die to individual leads of the lead frame, and a hard plastic encapsulant material which covers the other components and forms the exterior of the package.\nThe lead frame is the central supporting structure of such packages. A portion of the lead frame is internal to the package, i.e., completely surrounded by the plastic encapsulant. Portions of the leads of the lead frame extend externally from the package or are partially exposed within the encapsulant material for use in electrically connecting the die to the underlying substrate. Due to the relatively short electrical path from the die to the lead, such MLF packages can exhibit improved electrical efficiency in comparison to conventional Quad Flat Package (“QFP”) designs.\nOne of the drawbacks associated with certain MLF semiconductor packages is the limited surface area available to accommodate leads. Modern integrated circuit dies often include larger numbers of input and output pads than previous dies. To accommodate these additional pads, increasing numbers of leads must also be provided with semiconductor packages. Unfortunately, conventional MLF designs typically require these leads to be located on the bottom surface of the semiconductor package around the perimeter of a die paddle portion of the lead frame. Thus, the limited surface area available for lead placement on such packages can significantly impact the number of leads provided.\nThe present invention addresses this problem by providing various chip on lead configurations (“COL”) which position portions of leads underneath the integrated circuit die, thereby allowing more efficient use of the semiconductor package bottom surface area."}
{"text": "Since the introduction of the personal computer, in addition to the typical desktop supplies such as pencils, pens, sticky notes and paper clips, items associated with the personal computer, such as diskettes, CD ROM disks and CD marking pens, are required to be readily available, and preferably close at hand to the desktop. Storage of such items upon the desktop reduces the amount of desktop space available, which is compounded by the reduction in desktop space occasioned by the placement of the computer keyboard and monitor upon the desktop. Various desktop organizers, for holding pens, pencils, paper clips, sticky notes, rubber bands and the like have been known prior to the introduction of personal computers, and are used and directed toward the optimization of useable desktop space. Further, since the introduction of the computer, various desktop organizers and storage devices have been devised which are attached to the computer's components and are directed toward organizing supplies and optimizing usable desktop space where a computer is involved. However, these prior art devices which provide storage of desktop supplies and sundry computer related items, are separate component accessories and still occupy precious desktop space.\nIn an effort to address the shortcomings of the prior art devices, an integral keyboard and storage organizer of U.S. patent application Ser. No. 10/053,747 was created by DeVolpi. Since the creation of the device of U.S. patent application Ser. No. 10/053,747, further improvements have been made that enhance its utility including the addition of, latching means, various locking means, and configurations allowing access to electronic components and batteries incident to wireless embodiments.\nAccordingly, it is an object of the present invention to provide an improved integral keyboard and storage organizer which can store desktop supplies, sundry computer related items, and personal items, such as keys without necessitating the use of any additional space near the keyboard and incorporate latching means, various locking means, and configurations allowing access to electronic components and batteries incident to wireless embodiments."}
{"text": "1. Field\nMethods and apparatuses consistent with exemplary embodiments relate to estimating a malignant tumor in an object, and more particularly, to analyzing medical images of an object to estimate a malignant tumor.\n2. Description of the Related Art\nA probability that a woman develops breast cancer at some point in her life is more than 8%. At present, the most effective method of diagnosing breast cancer at an early stage is by a mammography image.\nA mammography image may include a lump-like mass, and a medical professional may detect the mass and determine whether the detected mass is benign or malignant.\nThe accuracy of mammography image analysis depends on the experience and knowledge of the medical professional, and an accuracy of 65% to 75% is possible.\nA device for analyzing a mammography image to determine whether a malignant tumor exists in the mammography image may use computer-aided diagnosis (CADx), and when a medical professional performs analysis using a CADx system, accuracy of the analysis may increase. However, the related art CADx system may miss detecting a feature suitable for determining a presence of malignant tumor.\nThus, an efficient CADx system is needed to increase the accuracy of a mammography image analysis."}
{"text": "A gate valve is a valve that opens by lifting a round or rectangular gate out of the path of the fluid. The gate valve has sealing surfaces between the gate and seats that are planar, leading to gate valves being used when a straight-line flow of fluid and minimum restriction is desired. Generally, gate valves are used to permit or prevent the flow of liquids, for example in a pipe line system.\nLiquid is able to flow when the gate of the gate valve is linearly retracted to open the flow path, whereas the liquid is prevented from flowing when the gate is linearly extended to a closed position. In the closed position liquid impinges a front side face of the gate and is prevented from passing between the gate and a body of the gate valve by rubber valve liners that surround the gate. Typically, the rubber valve liners are fixed in the body of the gate valve and rub against the gate as the gate opens and closes causing the rubber valve liners to wear away and lose their sealing ability. Furthermore, increasing compression of the rubber valve liners against the gate increases the wear rate of the rubber valve liners. Some gate valves employ spring components to maintain compression of the rubber valve liners, but compression force decreases as the rubber valve liners wear."}
{"text": "This invention relates to the hot rolling of elongated product, such as steel bars. More specifically, the invention concerns means and methods of in-line heat treating and finishing of the product.\nVarious systems are known for heat treatment of hot rolled products, such as bars. However, when the billets or blooms are produced in a caster such processes and facilities normally utilize off-line heat treating of the cast and hot rolled stock. It is known to perform certain limited heat treatments on the cast product before rolling. For example, U.S. Pat. No. 5,634,512 describes in-line thermal surface treatment of continuously cast blooms using cooling sprays to produce a bloom surface temperature of about 400-900xc2x0 C. after tempering caused by the hot core of the bloom. U.S. Pat. No. 4,786,338 discloses a system wherein product from a finishing mill is followed by a cooling line made up of a quenching zone and a temperature recovery zone before being passed to cooling beds for obtaining low temperature toughness. Japanese Published Patent Application No. Hei 63-149316 describes a method of producing steel bars and wire rods through continuous hot rolling for obtaining a directly spheriodized structure in which spheriodizing annealing is facilitated or possibly even eliminated.\nSuch prior art has the drawback of being limited in the capability of producing a wide variety of metallurgical structures in the rolled stock. Off-line heat-treatments to produce various microstructures normally require treatment times of 10-15 hours.\nThe apparatus of this invention is disposed in an intermediate position between a hot-rolling mill and a cold-finishing area for cleaning the rolled stock, cutting it to final form, bundling and packaging it for in-line production of a final product. The apparatus includes, in-line, and downstream of a hot rolling mill train, a thermocontrolled rolling device including a controlling temperature device and a sizing mill, shears, a quenching box, optionally an induction heater, a layers preparation system where layers of cut stock, such as bars, are consolidated prior to transfer to an annealing furnace, and thereafter to a cooling bed.\nBy selection of particular ones of these apparatus elements for use in particular cases of heat treatment, and by selection of the conditions of heat treatment, a great variety of product structures can be achieved, for example, a spheroidized steel product, a shearability or workability annealed product, a recrystallation annealed product, a solution annealed product, a quenched product (martensitic or bainitic quenching), or a quenched and tempered (quenched and stress-relieved) product.\nWith practice of the in-line heat treatments of this invention, prior art off-line treatment times of 10-15 hours can be reduced to about 1 to 2 hours.\nFor a better understanding of the invention, its operating advantages and the specific objectives obtained by its use, reference should be made to the accompanying drawings and description which relate to preferred embodiments thereof."}
{"text": "1. Field of the Invention\nThe present invention relates to a parallel plate plasma CVD (Chemical Vapor Deposition) apparatus, and in particular, to a parallel plate plasma CVD apparatus used in production of a titanium oxide thin film, which is used for a microwave monolithic IC (hereinafter referred to as MMIC).\n2. Description of the Related Art\nVarious methods for forming metal oxide thin films or dielectric thin films have been proposed. Examples of such methods are disclosed in, for example, Japanese Patent Publication (Kokoku) Nos. 59-37566, 60-12773, and 7-25545. Included in these publications is a report that a titanium oxide thin film can be formed using a CVD process. However, to carry out this method the substrate temperature must be high, for example, 500.degree. C. or more, during the film deposition step. This adversely affects the underlying substrate used for the device formation as well as other devices formed on the substrate. For example, if the temperature is high when MMIC capacitors are formed, characteristics of the underlying substrate such as GaAs and transistors deteriorate. To reduce this problem, a film deposition method which uses plasma generated in a deposition chamber during a CVD process (herein after referred to as a plasma CVD process) has been frequently used.\nIn a conventional parallel plate plasma CVD apparatus using a radio-frequency generator (13.56 MHz), a gaseous mixture is supplied into a glow discharge region between a top electrode (a radio-frequency input electrode) and a bottom electrode (a ground electrode), and activated for reaction in the mixture to deposit a film on a substrate located on the bottom electrode. The plasma must have a considerably high power density to decompose the gases sufficiently.\nHowever, since the surface of the substrate on which the film is deposited is located near the plasma generation region, high-energy particles generated by plasma discharge frequently attack the substrate surface. This may cause deterioration of film characteristics, that is, the film may be damaged or deposited abnormally by the impact of particles having excessive energy.\nIn order to solve the problem, proposals have been made to moderate the particle impact by, for example, controlling the self-bias near the substrate. For example, J. W. Coburn et al. disclosed a method for controlling ions and radicals incident on the substrate by applying a direct current flow without grounding the bottom electrode supporting the substrate in J. Appl. Phys., 43, 4965 (1972). A. Matsuda et al. disclosed a triode configuration consisting of parallel plate electrodes and a mesh electrode interposed therebetween in Thin solid Films, 92, 171 (1982).\nFIG. 3 is a schematic view illustrating the fundamental components of a conventional triode-type plasma CVD apparatus. A bottom electrode 2 and a top electrode 3 are arranged opposite to one other in a reaction chamber (not shown). The bottom electrode 2 is connected to an external DC electrical power source 4. A substrate S on which a dielectric thin film is to be formed is placed on the bottom electrode 2. The top electrode 3 is connected to a radio-frequency generator 5. A grooved mesh electrode 6 is located between the bottom electrode 2 and the top electrode 3 and a DC bias current is applied between the mesh electrode 6 and the bottom electrode 2.\nIn all of the above-mentioned methods for moderating particle impact, the bottom electrode must float electrically. However, as shown in FIG. 3, the bottom electrode of the conventional plasma CVD apparatus is generally connected to the body of the reaction chamber and grounded. Thus, the conventional plasma CVD apparatus cannot be used to carry out the foregoing methods.\nAccordingly, there is a need for a parallel plate plasma CVD apparatus which moderates damage to the substrate during the film deposition process, forms high-quality dielectric thin films, and permits the use of a modified conventional plasma CVD apparatus."}
{"text": "The present invention relates generally to a sludge treatment system for treating sludge, such as sewage sludge received from a municipality, and more particularly to an improved sludge treatment system which produces a pelletized fertilizer or soil amendment product from the sewage sludge.\nOther systems have been used to produce a pelletized product from sewer sludge. For example, U.S. Pat. No. 4,872,998 to Dausman et al. discloses an apparatus and process for forming a uniform pelletizable sludge product. The sludge is received by a holding tank which agitates the sludge to keep the solids in suspension with the liquid. The sludge is released from the holding tank through a control valve for mixing with flocculation and coagulation promoting chemicals and is then transported to a belt-type dewatering press. After the dewatering step, the sludge has 10-45% solids and is transported to a temporary storage tank. The sludge is fed from the temporary storage tank to an indirect dryer which is heated by either steam or hot oil from a boiler. Water vapor and other volatiles are driven off during the drying stage and pass from the dryer into a cyclone which removes the airborne dust. The resultant air and vapors are then pulled through a water jet scrubber to remove any remaining particles. This complete removal of all particulate matter is understood to be critical to subsequent process steps, and is believed to be a significant disadvantage of this process. The remaining particulate-free exhaust air is then passed to an odor control system. Alternatively, this remaining exhaust air may be returned to the boiler where the combustible components are ignited with the fuel used to heat the boiler. The dried sludge is fed from the indirect dryer to a storage tank or to a pelletizer. The pelletized material is then fed to a pellet storage tank. Thus, the Dausman et al. system processes the sludge without recycling of any previously heat-processed sludge.\nU.S. Pat. No. 3,963,471 to Hampton discloses a method for pelletizing sewage sludge for use as a fertilizer. Sludge is dewatered by adding coagulation and flocculation chemicals thereto and then passing the mixture through a dewatering press to form sludge cake. The dewatered sludge cake is mixed with a portion of the pellets constituting the final product which are in a hot, dried state. This mixing step lowers the heat required in the dryer to produce the final product and thereby prevents formation of clinkers and ignition of the organic material in the dryer. The larger the percentage of final product returned to the mixing chamber, the smaller in size are the resultant pellets formed in the dryer.\nU.S. Pat. No. 2,977,214 to McLellan discloses a method for pelletizing dried sewage sludge. The dried sludge is compacted between two compacting rollers into a rigid highly frangible sheet. A rotating breaker fragments the sheet into gross flakes or chunks. The Chunks are gravity fed from the compactor to a granulator. The granulator includes two sets of corrugated granulating rollers which grind the chunks received from the compactor. A screening mechanism sorts the granulated sludge into undersized, oversized and acceptable sized granules. The undersized granules are recycled to the compactor, and the oversized granules are recycled to the granulator.\nThe reissued U.S. Pat. No. RE.31,185 to Greenfield et al. discloses a process for dehydrating sewage sludges in a multistage evaporator. Oil is added to the sludge in the evaporator to yield a liquid pumpable mixture, which remains pumpable even after water removal. Vapors are removed at each stage of the multistage evaporator and used as a combustion source for the preceding step in the process. A portion of the anhydrous slurry containing sludge and the oil removed from the last stage is recycled back to a position upstream of the first stage, thereby controlling the viscosity of the evaporator feed. After the final evaporator stage, the resultant dry solids and oil are passed through a centrifuge for separation, and the separated oil is recycled. The solids from the centrifuge may be used for fertilizer or to fuel the boiler-furnace which supplies steam to the multistage evaporator.\nU.S. Pat. No. 4,321,151 to McMullen discloses a process for converting sewage sludge into energy useful for operating a sewage treatment plant. Raw settled sludge is withdrawn from the treatment stream early in the process sequence to preserve the energy content of the sludge. The withdrawn sludge is dewatered, caked and fragmented for pyrolysis to produce a combustible gas which may be used to fuel a gas turbine-generator to power the plant. The combustible gas is cleaned by removing the suspended tars and oils therein. The dried feed stock is pneumatically conveyed for pyrolyzing to a combustion boiler-pyrolyzer, and a portion is diverted to a rotary mixing device or blender. In the blender, the dried feed stock is blended with residual tar, oil and char (pyrolysis by-products) produced from pyrolysis of mixed and dried feed stock. The blended stock is directed to a pelletizing unit to produce pellets one inch in diameter and one and one-half inches long. The pellets are pneumatically conveyed to a combustion zone of the boiler-pyrolyzer.\nU.S. Pat. No. 3,960,725 to Bjermo et al. discloses a method of converting heavy-metal hydroxides into heavy-metal oxide powders or granules using concentration sintering. Before sintering, the sludge is dewatered and deposited on a media carrier in a 1-5 mm thick film and the film is then dried. This method is useful for eliminating the toxicity of the metal hydroxy sludge for environmental considerations.\nU.S. Pat. No. 4,097,378 to St. Clair discloses a method for removing water from peat and various sludges. The sludge is mixed with oil and fed into a condenser where it is heated by direct contact with steam. Liberated uncondensed gases are routed to a conventional vent system.\nU.S. Pat. No. 4,402,834 to Bastgen et al. discloses a continuous dewatering process for sludges principally involving filtration. The raw sludge is mixed with a flocculating agent in a rotary filter drum. The thickened sludge is conveyed through a filtering assembly and subjected to dewatering forces including vacuum, centrifugal, static pressure or pressure while moving along the filter surfaces.\nU.S. Pat. No. 4,795,568 to Chen discloses a waste water treatment system using excessive amounts of pressurized oxygen gas to oxidize the waste water solution and increase the evaporation of waste therefrom. The oxidized liquid effluent undergoes a series of flashing steps to produce a saturated solution Which crystallizes upon cooling.\nU.S. Pat. No. 3,342,731 to Baumann discloses a method of dewatering fresh or digested sludges or concentrate resulting from centrifuging fresh or digested sludge. After centrifuging, the sludge is further dewatered by filtering and the dewatered sludge is then incinerated. The ash resulting from the incineration of the dewatered sludge is added to the digested sludge or fresh sludge, along with lime for improved filterability. The lime remaining after incineration of sludge that has been treated with the ash and lime may be recycled by mixing it with the fresh lime and ash to be added to the sludge.\nU.S. Pat. No. 1,915,240 to Putnam discloses a sewage purification system where lime and ferric chloride are mixed with the sludge to kill bacteria, oxidize and deodorize the treated mass and to promote flocculation. The sludge is centrifuged or-otherwise dewatered then dried in a sludge dryer and ground for delivery to a storage bin. The resulting dried ground sludge may be used as a fertilizer or burned as fuel in the dryer. Alternatively, instead of dewatering, drying and grinding the sludge for fertilizing or heat, the sludge may be passed directly from a dewaterer and processed through a retort to produce char or carbonized particles. This resulting char product is introduced into the raw sewage before the chemical treatment step to absorb the gases and putrescible matter and increase the speed at which the solids settle from the liquid.\nU.S. Pat. No. 3,025,151 to Berg discloses a method for digesting solids separated from sewage sludge. The separated solids are mixed with a biological culture and aerated at a controlled temperature to decompose the fibrous cellulose, reduce the grease content of the sewage and remove objectionable odors from the solids. The resulting product may be pelletized for use as a fertilizer.\nU.S. Pat. No. 3,676,074 to Shibayama et al. discloses a cylindrical rotary tank having an air pipe along the axis thereof. Organic materials are fed into the cylindrical tank and fermented using aerobic bacteria and thereby causing substantially no odor nuisance. The fermentation process also heats the tank to aid in drying the organic waste and produce a solid granular material which may be used as a fertilizer.\nU.S. Pat. Nos. 4,660,628 to Solberg et al. and 3,800,865 to Onarheim et al. each disclose rotary heat exchangers wherein the heating or cooling medium is passed through a hollow shaft having a plurality of fins extending outwardly therefrom. The shaft is enclosed within a container which receives the moist, solid or semi-solid materials which are to be heated or cooled.\nThe Japanese Patent No. 5,439,372 illustrates the method of adding aluminum sulfate to waste sludge for settling, and thereafter adding a macromolecule type aggregating agent. The mixture is then centrifugally separated and the solid matter heated for solidification in a drying furnace at 200-500% C to obtain a dried cake.\nU.S. Pat. No. 3,695,432 to McCrink discloses a sewage disposal system which pumps solid wastes separated from sewage through a heating retort. The solids are mixed with a non-oxidizing gas such as ammonia, methane or propane which is injected therein under pressure, and the mixture is heated to produce a fine powdery ash. The gases produced during the process are passed through a flame curtain prior to exiting the system to remove all odors therefrom.\nU.S. Pat. No. 548,561 to Lamb discloses an apparatus for making fertilizer from refuse. The refuse is filtered through sand and then evaporated using steam passing through a coil and then further dried using steam or hot air to produce a fertilizer. Odors produced by this system are conducted to the furnace for destruction therein.\nU.S. Pat. No. 659,503 to Wood discloses a sewage treating apparatus which receives sludge solids in a vibrating and oscillating hopper comprising a filter-bed holder. Water is removed by the vibrating and oscillating hopper, and the resulting solids entrapped therein pass to a drying cylinder. F where the material is thoroughly dried and granulated by heating and revolving the drying cylinder to produce a fertilizer.\nU.S. Pat. No. 3,909,410 to Neukamm discloses a sewage sludge treating process wherein the sludge slurry is continuously distributed in a rotating contactor drum onto preheated aggregate pieces which are continually fed into the drum. A portion of the slurry instantaneously adheres to the heated aggregate and the non-adhering slurry exits the drum for recycling. The coated aggregate is collected from the drum and heated during travel to a tumbler drum. In the tumbler drum, the sludge coating is broken away from the aggregate and pulverized into a powder or dust which is the product. The decoated aggregate is recycled into the rotating contactor drum.\nFurthermore, raw primary or a mixture of raw primary and activated sewage slush solids contain a significant portion of fibrous materials which is derived from paper products or other substances. When the sludge solids are dried, the fibrous material becomes visible and undesirable and is referred to as \"fluff\", because it tends to expand greatly in volume Prior art drying systems typically remove the fluff and dispose of it as a separate waste product, or process the fluff off-site. If raw primary sludge or a mixture of raw primary and activated sludge solids is digested in an aerobic or anaerobic digestion process, a large portion of the paper or cellulose is converted and the fluff is substantially reduced.\nNone of the above sewage treatment systems is understood to process the sewage sludge into a pelletized fertilizer having a long shelf life, such as on the order of seven months. Furthermore, none of these known systems is understood to make use of waste oil, such as transformer oil, crankcase oil, and other reclaimed oil in a heating stage of the sewage treatment process without venting significantly detectable amounts of gaseous pollutants thereby produced to atmosphere. Furthermore, none of the above sewage treatment systems is understood to fulfill such deficiencies with a minimal odor dispersing apparatus which may be used in a urban setting.\nThus, a need exists for an improved sludge treatment system for pelletizing sewage sludge, including that which contains significant amounts of fibrous materials, such as may be used to produce a pelletized fertilizer or soil amendment product, which is not susceptible to the above limitations and disadvantages."}
{"text": "A spiral inductor shows a skin effect where, as frequency is increased, current is concentrated on a surface of a conductor. A spiral inductor also shows a proximity effect where upon bringing another conductor in proximity with the spiral inductor, a high frequency current is more intensively flowed into a portion close to the other conductor. In order to represent these two effects in an equivalent circuit of the inductor, a coupling form of a 2-π equivalent circuit model that is a modification of an existing single-π model was recently disclosed in an article and is shown in FIG. 1.\nHowever, the 2-π equivalent circuit model has problems in that the extraction of the model is complicated and a large error can be generated in the value of the extracted parameter."}
{"text": "1. Field of the Invention\nThe invention relates to N-substituted benzene sulfonamides which inhibit β-amyloid peptide release and/or its synthesis and, therefore, are useful in the prevention of cognitive disorders in patients susceptible to cognitive disorders and/or in the treatment of patients with cognitive disorders in order to inhibit further deterioration in their condition.\n2. State of the Art\nAlzheimer's Disease (AD) is a degenerative brain disorder characterized clinically by progressive loss of memory, cognition, reasoning, judgment and emotional stability that gradually leads to profound mental deterioration and ultimately death. AD is a very common cause of progressive mental failure (dementia) in aged humans and is believed to represent the fourth most common medical cause of death in the United States. AD has been observed in races and ethnic groups worldwide and presents a major present and future public health problem. The disease is currently estimated to affect about two to three million individuals in the United States alone. AD is at present incurable. No treatment that effectively prevents AD or reverses its symptoms and course is currently known.\nThe brains of individuals with AD exhibit characteristic lesions termed senile (or amyloid) plaques, amyloid angiopathy (amyloid deposits in blood vessels) and neurofibrillary tangles. Large numbers of these lesions, particularly amyloid plaques and neurofibrillary tangles, are generally found in several areas of the human brain important for memory and cognitive function in patients with AD. Smaller numbers of these lesions in a more restrictive anatomical distribution are also found in the brains of most aged humans who do not have clinical AD. Amyloid plaques and amyloid angiopathy also characterize the brains of individuals with Trisomy 21 (Down's Syndrome) and Hereditary Cerebral Hemorrhage with Amyloidosis of the Dutch Type (HCHWA-D). At present, a definitive diagnosis of AD usually requires observing the aforementioned lesions in the brain tissue of patients who have died with the disease or, rarely, in small biopsied samples of brain tissue taken during an invasive neurosurgical procedure.\nThe principal chemical constituent of the amyloid plaques and vascular amyloid deposits (amyloid angiopathy) characteristic of AD and the other disorders mentioned above is an approximately 4.2 kilodalton (kD) protein of about 39–43 amino acids designated the β-amyloid peptide (BAP) or sometimes Aβ, AβP or β/A4. β-Amyloid peptide was first purified and a partial amino acid sequence was provided by Glenner et al., Biochem. Biophys. Res. Commun., 120: 885–890 (1984) The isolation procedure and the sequence data for the first 28 amino acids are described in U.S. Pat. No. 4,666,829.\nMolecular biological and protein chemical analyses have shown that the β-amyloid peptide is a small fragment of a much larger precursor protein termed the amyloid precursor protein (APP), that is normally produced by cells in many tissues of various animals, including humans. Knowledge of the structure of the gene encoding APP has demonstrated that β-amyloid peptide arises as a peptide fragment that is cleaved from APP by protease enzyme(s). Sequential processing of the precursor protein by the enzymes referred to generically as beta- and gamma-secretases, give rise to the β-amyloid peptide fragment. Both enzymes have now been molecularly cloned, and characterized to differing levels.\nSeveral lines of evidence indicate that progressive cerebral deposition of β-amyloid peptide plays a seminal role in the pathogenesis of AD and can precede cognitive symptoms by years or decades. See, for example, Selkoe, Neuron, 6:487–498 (1991). The most important line of evidence is the discovery that missense DNA mutations at amino acid 717 of the 770-amino acid isoform of APP can be found in affected members but not unaffected members of several families with a genetically determined (familial) form of AD (Goate et al., Nature, 349:704–706 (1990); Chartier Harlan et al., Nature, 353:844–846 (1989); and Murrell et al., Science, 254: 97–99 (1991.) Another such mutation, known as the Swedish variant, is comprised of a double mutation changing lysine595-methionine596 to asparagine595-leucine596 (with reference to the 695 isoform was found in a Swedish family) was reported in 1992 (Mullan et al., Nature Genet., 1:345–347 (1992). Genetic linkage analyses have demonstrated that these mutations, as well as certain other mutations in the APP gene, are the specific molecular cause of AD in the affected members of such families. In addition, a mutation at amino acid 693 of the 770-amino acid isoform of APP has been identified as the cause of the β-amyloid peptide deposition disease, HCHWA-D, and a change from alanine to glycine at amino acid 692 appears to cause a phenotype that resembles AD is some patients but HCHWA-D in others. The discovery of these and other mutations in APP in genetically based cases of AD prove that alteration of APP metabolism, and subsequent deposition of its β-amyloid peptide fragment, can cause AD.\nDespite the progress which has been made in understanding the underlying mechanisms of AD and other β-amyloid peptide related diseases, there remains a need to develop methods and compositions for treatment of the disease(s). Ideally, the treatment methods would advantageously be based on drugs which are capable of inhibiting β-amyloid peptide release and/or its synthesis in vivo.\nOne approach toward inhibiting amyloid peptide synthesis in vivo is by inhibiting gamma secretase, the enzyme responsible for the carboxy-terminal cleavage resulting in production of β-amyloid peptide fragments of 40 or 42 residues in length. The immediate substrates for gamma secretase are β-cleaved, as well as α-cleaved carboxy-terminal fragments (CTF) of APP. The gamma-secretase cleavage site on β- and α-CTF fragments occurs in the predicted transmembrane domain of APP. Inhibitors of gamma-secretase have been demonstrated to effect amyloid pathology in transgenic mouse models (Dovey, H. F., V. John, J. P. Anderson, L. Z. Chen, P. de Saint Andrieu, L. Y. Fang, S. B. Freedman, B. Folmer, E. Goldbach, E. J. Holsztynska et al. (2001). “Functional gamma-secretase inhibitors reduce beta-amyloid peptide levels in brain.” J Neurochem 76(1): 173–81.)\nGamma secretase is recognized to be a multi-subunit complex comprised of the presenilins (PS1 or PS2), Nicastrin, Aph-1, and Pen 2 (De Strooper, B. (2003). “Aph-1, Pen-2, and Nicastrin with Presenilin generate an active gamma-Secretase complex.” Neuron 38(1): 9–12; Edbauer, D., E. Winkler, J. T. Regula, B. Pesold, H. Steiner and C. Haass (2003). “Reconstitution of gamma-secretase activity.” Nat Cell Biol 5(5): 486–8; Kimberly, W. T., M. J. LaVoie, B. L. Ostaszewski, W. Ye, M. S. Wolfe and D. J. Selkoe (2003). “Gamma-secretase is a membrane protein complex comprised of presenilin, nicastrin, Aph-1, and Pen-2.” Proc Natl Acad Sci USA 100(11): 6382–7). Much evidence indicates that PS comprises the catalytic moiety of the complex, while the other identified subunits are necessary for proper maturation and sub-cellular localization of the active enzyme complex (reviewed in De Strooper, B. (2003). “Aph-1, Pen-2, and Nicastrin with Presenilin generate an active gamma-Secretase complex.” Neuron 38(1): 9–12.) Consistent with this hypothesis: PS knock-out mice exhibit significant reductions in β-amyloid production (De Strooper, B., P. Saftig, K. Craessaerts, H. Vanderstichele, G. Guhde, W. Annaert, K. Von Figura and F. Van Leuven (1998). “Deficiency of presenilin-1 inhibits the normal cleavage of amyloid precursor protein.” Nature 391(6665): 387–90; Haass, C. and D. J. Selkoe (1998). “Alzheimer's disease. A technical KO of amyloid-beta peptide.” Nature 391(6665): 339–40; Herreman, A., L. Serneels, W. Annaert, D. Collen, L. Schoonjans and B. De Strooper (2000). “Total inactivation of gamma-secretase activity in presenilin-deficient embryonic stem cells.” Nat Cell Biol 2(7): 461–2); point mutations of putative active site aspartate residues in PS trans-membrane domains inhibit β-amyloid production in cells in a dominant negative fashion (Wolfe, M. S., W. Xia, B. L. Ostaszewski, T. S. Diehl, W. T. Kimberly and D. J. Selkoe (1999). “Two transmembrane aspartates in presenilin-1 required for presenilin endoproteolysis and gamma-secretase activity.” Nature 398(6727): 513–7; Kimberly, W. T., W. Xia, T. Rahmati, M. S. Wolfe and D. J. Selkoe (2000). “The transmembrane aspartates in presenilin 1 and 2 are obligatory for gamma-secretase activity and amyloid beta-protein generation.” J Biol Chem 275(5): 3173–8); active site directed substrate-based transition state isosteres designed to inhibit gamma secretase directly conjugate to PS (Esler, W. P., W. T. Kimberly, B. L. Ostaszewski, T. S. Diehl, C. L. Moore, J. Y. Tsai, T. Rahmati, W. Xia, D. J. Selkoe and M. S. Wolfe (2000). “Transition-state analogue inhibitors of gamma-secretase bind directly to presenilin-1.” Nat Cell Biol 2(7): 428–34; Li, Y. M., M. Xu, M. T. Lai, Q. Huang, J. L. Castro, J. DiMuzio-Mower, T. Harrison, C. Lellis, A. Nadin, J. G. Neduvelil et al. (2000). “Photoactivated gamma-secretase inhibitors directed to the active site covalently label presenilin 1.” Nature 405(6787): 689–94); finally, allosteric gamma secretase inhibitors have likewise been demonstrated to bind directly to PS (Seiffert, D., J. D. Bradley, C. M. Rominger, D. H. Rominger, F. Yang, J. E. Meredith, Jr., Q. Wang, A. H. Roach, L. A. Thompson, S. M. Spitz et al. (2000). “Presenilin-1 and -2 are molecular targets for gamma-secretase inhibitors.” J Biol Chem 275(44): 34086–91.)\nCurrent evidence indicates that in addition to APP processing leading to β-amyloid synthesis, gamma-secretase also mediates the intra-membrane cleavage of other type I transmembrane proteins (reviewed in Fortini, M. E. (2002). “Gamma-secretase-mediated proteolysis in cell-surface-receptor signalling.” Nat Rev Mol Cell Biol 3(9): 673–84, see also Struhl, G. and A. Adachi (2000). “Requirements for presenilin-dependent cleavage of notch and other transmembrane proteins.” Mol Cell 6(3): 625–36.) Noteworthy among the known substrates of gamma-secretase is mammalian Notch 1. The Notch 1 protein is important for cell fate determination during development, and tissue homeostasis in the adult. Upon ligand engagement via the Notch ecto-domain, Notch undergoes sequential extra-cellular and intra-membrane processing analogous to APP. The intra-membrane processing of Notch mediated by gamma secretase leads to release of the Notch intracellular domain (NICD). The NICD fragment mediates Notch signaling via translocation to the nucleus, where it regulates expression of genes mediating cellular differentiation in many tissues during development, as well as in the adult.\nDisruption of Notch signaling via genetic knock-out (KO) results in embryonic lethal phenotype in mice (Swiatek, P. J., C. E. Lindsell, F. F. del Amo, G. Weinmaster and T. Gridley (1994). “Notchl is essential for postimplantation development in mice.” Genes Dev 8(6): 707–19; Conlon, R. A., A. G. Reaume and J. Rossant (1995). “Notchl is required for the coordinate segmentation of somites.” Development 121(5): 1533–45.) The Notch KO phenotype is very similar to the phenotype observed PS1 KO mice, and precisely reproduced by PS1/PS2 double KO mice (De Strooper et al. (1998). “Deficiency of presenilin-1 inhibits the normal cleavage of amyloid precursor protein.” Nature 391(6665): 387–90; Donoviel, D. B., A. K. Hadjantonakis, M. Ikeda, H. Zheng, P. S. Hyslop and A. Bernstein (1999). “Mice lacking both presenilin genes exhibit early embryonic patterning defects.” Genes Dev 13(21): 2801–10; Herreman, A., L. Serneels, W. Annaert, D. Collen, L. Schoonjans and B. De Strooper (2000). “Total inactivation of gamma-secretase activity in presenilin-deficient embryonic stem cells.” Nat Cell Biol 2(7): 461–2.) This convergence of phenotypes observed in knock-out mice of either the substrate (Notch) or the enzyme (PS) suggests that inhibitors of gamma secretase that also inhibit Notch function may be limited as therapeutic agents owing to the importance of Notch function in adult tissues (Fortini, M. E. (2002). “Gamma-secretase-mediated proteolysis in cell-surface-receptor signalling.” Nat Rev Mol Cell Biol 3(9): 673–84.) As APP knock-out mice develop normally and without an overt phenotype Zheng, H., M. Jiang, M. E. Trumbauer, R. Hopkins, D. J. Sirinathsinghji, K. A. Stevens, M. W. Conner, H. H. Slunt, S. S. Sisodia, H. Y. Chen et al. (1996). “Mice deficient for the amyloid precursor protein gene.” Ann N Y Acad Sci 777:421–6; Zheng, H., M. Jiang, M. E. Trumbauer, D. J. Sirinathsinghji, R. Hopkins, D. W. Smith, R. P. Heavens, G. R. Dawson, S. Boyce, M. W. Conner et al. (1995). “beta-Amyloid precursor protein-deficient mice show reactive gliosis and decreased locomotor activity.” Cell 81(4): 525–31, the cumulative evidence, therefore, suggests that preferred gamma secretase inhibitors would have selectivity for inhibiting gamma secretase processing of APP over gamma secretase processing of Notch."}
{"text": "Actuators typically are mechanical devices that are used for moving or controlling a mechanism, system or the like and typically convert energy into some type of motion. Examples of actuators can be found in any number of applications encountered in everyday life including automotive, aviation, construction, farming, factories, robots, health care and prosthetics, among other areas.\nMobile robotics and advanced prosthetics will likely play important roles in the future of the human race. Actuators frequently are used in these applications that enable movement of a robot or user arm or other appendage or item as desired.\nMost existing mobile robots and advanced prosthetics, however, lack the strength and speed necessary to be effective. This is because they suffer from poor specific power (strength×speed/weight) which determines how quickly work can be done compared to another actuator of the same weight.\nFor example, if such devices are capable of lifting significant weight, they must do so very slowly, which inhibits their adoption for most applications. On the other hand, devices that can move more quickly are just not capable of handling anything more than the smallest weight."}
{"text": "1. Field of the Invention\nThe present invention relates to an article for reducing or eliminating parasitic oscillation (PL) and amplified spontaneous emission (ASE) in solid state lasers. More specifically, it relates to a thin film attached to strategic facets of a solid state laser gain material to frustrate the total internal reflection of parasitic oscillation and amplified stimulated emission.\n2. Description of Related Art\nThe architecture of most high average power solid state lasers consists of a laser medium interposed between a highly reflecting mirror and an output coupler. Solid state zigzag slab lasers are typically used when high average power is needed because the optical path utilizes more gain medium and tends to average out index of refraction gradients caused by the thermal loading. Diode lasers or flashlamps are often used to optically excite the gain medium. In slab lasers, the medium is optically excited through the two largest exposed planes which are parallel to the optical axis. The laser medium must be cooled due, in part, to the high optical fluxes coupled into the slab by the pumping mechanism. If not properly designed, the stored energy density in slab lasers decreases. This performance reduction is caused by the complementary effects of amplified spontaneous emission (ASE) and parasitic oscillations that occur in the slab laser. [see for example, \"Fluorescence Amplification and Parasitic Oscillation Limitations in Disk Lasers\", by J. B. Trenholme, NRL Memorandum Rep.2480, July 1972; J. E. Swain et al., J. Appl. Phys., Vol. 40 p.3973, (1969); and J. M. McMahon et al., IEEE J. Quantum Electron., QE-9, p.992, (1973); and U.S. Pat. No. 4,849,036].\nEdge claddings are used on solid state laser material to absorb the amplified spontaneous emission and to suppress the onset of parasitic oscillations that would otherwise reduce the stored energy. In general, these claddings consist of a material that is refractive index matched to the laser glass and which contains a dopant that absorbs at the laser (ASE) frequency. A number of different materials have been used for cladding, ranging from sprayed-on glass frits to liquids to castings of monolithic glass. [See, for example, G. Dube and N. L. Boling, Applied Optics, Vol. 13, p.699, (1974); G. Guch, Jr.,Applied Optics, Vol. 15, p.1453, (1976); and D. Milam et al., \"Platinum Particles in the Nd:doped Disks of Phosphate Glass in the Nova Laser\", Laser Induced Damage in Optical Materials 1985: Proceedings of the Boulder Damage Symposium, Boulder, Colo., November 1985] In the 100-kJ Nova pulsed laser system at Lawrence Livermore National Laboratory, claddings of monolithic glass doped with ionic copper that absorbs at 1 micrometer have been used. Although the performance of this latter cladding works well in high peak power lasers, the heat transfer properties are not adequate for use in high average power solid state lasers.\nParasitic lasing (PL) severely limits the performance of high gain solid state zigzag slab lasers by depleting energy in the upper laser level. The problem is enhanced because solid state laser materials have a higher index of refraction than the surrounding media. The traditional method for suppressing ASE and PL, as stated supra, has been to bind a relatively thick layer of weakly absorbing material with closely matched index of refraction to the edges. Index matching suppresses the reflection of light rays at the boundary between the two media, enabling the absorption of radiation within the bonding material. These traditional edge absorbers for ASE and PL suppression work adequately for high peak power, low average power lasers, but, due to unacceptable heat transfer properties, are not suitable for high average power lasers.\nIt is desirable to have a thin film attached onto the otherwise reflective surfaces of a solid state laser medium to bring about suppression of total internal reflections and reduction in the internal reflectivity for wide angular incidence without use of the doped index matching material. The present invention provides such an article."}
{"text": "Modern cotton harvesters typically include a relatively large high volume receptacle mounted to the frame of the harvester. Recently, the receptacle or basket assembly of cotton harvesters has been designed such that the height and, thus, the capacity of the basket assembly is adjustable. The receptacle or basket assembly on recent cotton harvester designs includes a lower basket portion mounted on the frame of the harvester with an upper basket portion being movably mounted and supported for elevational movement relative to the lower basket portion under the influence between an extended field working position and a retracted storage position.\nA conventional cotton harvester further includes a plurality of harvesting row units arranged across a front end of the harvester. Each harvesting unit has duct structure operably associated therewith for receiving cotton materials harvested by the respective harvesting unit and for delivering the cotton materials into the receptacle. To optimize receptacle capacity, the cotton conveying structure extends upwardly from a respective harvesting row unit to a top forward portion of the upper basket portion of the receptacle wherein one or more openings are provided in registry with an exhaust end of the cotton conveying duct structure. To promote the delivery of cotton from the duct structure into the receptacle, the upper basket portion of the receptacle can include hood structure that extends vertically over and in operable combination with the discharge end of the duct structure.\nAlthough beneficially increasing the capacity of the receptacle, maintaining the upper basket portion in an extended field working position presents certain problems and drawbacks. Because of shipping size limitations, the upper basket portion must be moved to its retracted storage position when the harvester is freighted by rail or truck. Because of power lines and the like, the increased height of the receptacle, with the upper basket portion in a raised field working position, makes it difficult to transport the harvester from one location to another. Moreover, the normal height of barn doors makes moving the harvester indoors for repairs or storage a problem because of the increased height added to the receptacle by the upper basket portion. Thus, and after the cotton materials are removed from the receptacle, the upper basket portion of the receptacle is lowered into its retracted position to facilitate transportation of the harvester.\nWhile the height of the duct structure can exceed the height of the cab region, it is beneficial for transportation purposes to lower the height of the duct structure along with reducing the height of the receptacle. Some cotton harvesters have a plurality of duct structures extending upwardly toward the opening in the receptacle from each harvesting row unit. As will be appreciated, removing portions of each duct structure to reduce the height thereof is a labor intensive manual process. Reassembling each duct structure to return the machine to an operational condition also involves a labor intensive effort. A customer, for example, who desires to transport his cotton harvester from one location to the other must often spend a considerable number of man hours in tearing down and setting up the receptacle, thereby reducing the overall productivity of the harvester.\nThus, there is a need and a desire for a large capacity receptacle for a cotton harvester having arranged in combination therewith duct structure designed to accommodate shipping and transportation of the harvester and which obviates the need for tearing down and setting up the entire duct structure in response to the position of the upper basket portion of the receptacle."}
{"text": "In drug discovery research, proper evaluation of incorporation, metabolism and excretion of a drug or the like in the hepatic tissue of a living body is extremely important. Further, for simply carrying out large-scale drug screening, or from an ethical point of view, an in vitro test method has been demanded. As a means for such a test, cultured cells are used, and their culture is preferably carried out under conditions reflecting those in a living body as much as possible. For example, as a culture system for hepatocytes, one which can excellently reproduce the proper polarity of the hepatic tissue of a living body, wherein, for example, the cell membrane is clearly differentiated between the blood vessel side and the bile canaliculus membrane side, is demanded\nFor example, Patent Document 1 describes that it is thought that, by binding a binding protein or an adhesion protein to a microscale substrate and culturing hepatocytes thereon, formation of a bile canaliculus is promoted. Further, Patent Document 2 describes that, by culturing small hepatocytes on a polycarbonate porous sheet covered with collagen, a bile canaliculus-like structure was formed on Day 30 of the culture.\nIn the above-described three-dimensional culture methods, preparation of hepatocytes having the polarity from the primary cultured cells takes several weeks to several months.\nOn the other hand, in cases where the collagen gel sandwich method (Non-patent Document 1: LeCluyse et al., Am J Physiol Cell Physiol, 1994, vol. 266, pp. 1764-1774) is used, formation of a bile canaliculus and the bile component excretion activity begins to be detected about 3 or 4 days after deposition of a collagen gel. However, even in the cases where the collagen gel sandwich method is used, several days are required before the bile component excretion activity is obtained. Further, metabolites cannot be continuously analyzed in such cases since a bile duct having an outlet, as is the case in a living body, is required for the continuous analysis. There is also a problem in that the number of bile canaliculi is too small and the activity is too low, to be used in drug screening.\nFurther, in Non-patent Document 2 (Tissue Engineering vol.13 Number 8 2007 2105-2117) and Non-patent Document 3 (Acta Biomaterialia 5 2009 613-620), it is shown that, when hepatocytes are to be cultured on gas-permeable PDMS (polydimethylsiloxane), adhesion of the cells can be maintained for 3 days by coating the surface of the PDMS with PEM or by forming thin pillars having a diameter of 1 to 3 μm. However, the adhesion cannot be maintained for a long time and hence those cultured cells are problematic in view of stable use in a test. Further, it has been generally thought that PDMS is not suitable for adherent culture of hepatocytes."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a method for RF output power control of a wireless communication device and, more particularly, to a method for output power control by switching the wireless communication device between different modes according to a state parameter corresponding to a received beacon.\n2. Description of the Prior Art\nIn recent years, applications in fields such as mobile communications as well as wireless local area network (WLAN) have attracted considerable attention. For wireless communication devices such as mobile phones, base stations, wireless network interface cards and access points, and so forth, it is necessary to reduce power consumption so as to prevent ageing and degradation of the components. Especially, portable wireless communication devices require batteries as power supply; therefore, it is necessary for such devices to reduce power consumption so as to prolong the battery lifetime. Moreover, in order to prevent human from absorbing the electromagnetic waves, it is required that the RF energy of the nuisance radiation from the wireless communication devices must be conserved when the devices are not operating or operating at a short distance.\nIn view of the power consumption, most of the power is consumed by the transmitter of wireless communication devices when the transmitter is transmitting RF signals. Accordingly, the present invention provides a method for RF output power control of a wireless communication device so as to reduce power consumption and eliminate nuisance radiation on one hand and leave the link quality of signal transmission unaffected on the other."}
{"text": "1. Technical Field\nEmbodiments relate to semiconductor devices. More particularly, embodiments relate to a semiconductor device and a floating prevention method thereof, which may prevent a portion of signal line from floating when a signal line is an open state.\n2. Description of the Related Art\nSemiconductor devices, e.g., DRAM, etc., generally employ a redundancy scheme for replacing a defective memory cell with a spare memory cell.\nFor efficiency, in the redundancy scheme, a replacement may be obtained in a unit of a column selection line or a normal word line enable signal. For example, when a defective column selection line is replaced with a spare column selection line, the defective column selection line is maintained at a ground level.\nHowever, e.g., when a line open state occurs in a defective column selection line, a portion of the defective column selection line may be in a floating state. Thus, electrical coupling with an adjacent column selection line may occur."}
{"text": "The present invention relates to a semiconductor memory device and a manufacturing method thereof. More particularly, the present invention relates to a ferroelectric memory device in which a ferroelectric film is used for a gate insulating film and a manufacturing method thereof.\nRecent advances in thin-film technology have produced nonvolatile memory devices using ferroelectric films (i.e., ferroelectric memory devices), which can perform high-speed read/write operations by using the polarization reversal and retention characteristics of the ferroelectric film. Since the polarization reversal of a ferroelectric film is caused by a transition at the atomic level, a ferroelectric memory device operates 10.sup.4 -10.sup.5 times faster than other types of nonvolatile memories (e.g., EEPROM or flash memory devices) and thus can achieve read-cycle speeds comparable to that of a DRAM device, i.e., in the hundreds of nanoseconds. Also, since polarization reversal only requires a low supply voltage, e.g., two to five volts, the memory device does not require the higher supply voltage, e.g., ten to twelve volts, necessary for EEPROM or flash memory read operations.\nBroadly speaking, there are two types of conventional ferroelectric memory devices: one which detects the amount of charge stored in a capacitor and one which detects a change in the resistance caused by the spontaneous polarization of a ferroelectric.\nThe first type comprises two transistors and two capacitors or of a single transistor-capacitor pair. Memory devices of this type are widely adopted in DRAM devices and are generally formed of CMOS transistors, with a thick interlayer insulating film separating the transistors from a ferroelectric capacitor formed on top. In this structure, although the impact of the ferroelectric electrode material on the underlying CMOS device can be reduced, a destructive read-out problem persists in that the stored data is ultimately destroyed by being read. This destructive read-out becomes of greater concern as the number of read/write operations in a memory device increases.\nThe second type of ferroelectric memory device endeavors to overcome this destructive read-out problem. An example of the second type of ferroelectric memory device is called a metal-ferroelectric insulator semiconductor field-effect transistor (MFIS FET). Such a device provides for nondestructive read-out, which makes it appropriate for increased read/write operations. In principle, this type of memory device, which has a single transistor compared to a DRAM having both a transistor and a capacitor, reduces the cell area without altering the cell structure in accordance with typical MOS FET scaling rules. Moreover, its read time need not be as long as that of a flash memory which is a nonvolatile memory, and data can be retained after a read-out operation. The memory device employs a metal-insulator-semiconductor (MIS) structure using a ferroelectric film as a gate insulating film, in which a read-out operation is performed by forming an inversion layer in a channel region of a transistor. The inversion layer is formed by controlling the potential on a silicon interface based on the polarization retention of a ferroelectric.\nHowever, there are additional problems encountered with the use of the MFIS FET. The major problem derives from the fact that in the MFIS FET a ferroelectric film is formed directly on a silicon semiconductor substrate. When an oxide ferroelectric film, such as PZT (PbZr.sub.x T.sub.1-x O.sub.3), is formed directly on a silicon substrate, oxygen atoms are injected from the oxygen-rich ferroelectric film into its interface with the silicon substrate, thereby resulting in the formation of a superfluous thin film, e.g., an SiO.sub.2 film. As a result, either the composition ratio of the ferroelectric film is locally destroyed, or a higher operating voltage is required. Furthermore, charged particles are injected into the film by trap levels resulting from stress of the ferroelectric film, thereby erasing the charge produced by polarization retention. Also, if the film is formed at a high temperature, constituents of the ferroelectric are apt to diffuse into the silicon substrate and alter the characteristics of the field effect transistor. In addition, if a non-oxide ferroelectric, e.g., BaMgF.sub.4, were used instead of an oxide ferroelectric, fluorine ions are apt to penetrate into the gate insulating film, resulting in the elimination of the polarization characteristics.\nAccordingly, since a ferroelectric film is not compatible with a silicon substrate in terms of lattice constant and thermal expansion coefficient, it is very difficult to form a high-quality ferroelectric film on the silicon substrate. Furthermore, to form the source/drain regions by self-alignment, a film showing tolerance to a thermal process at approximately 1,000.degree. C. would be required.\nIn an effort to solve the above problems, an MFIS FET in which a PbTiO.sub.3 film is deposited on a CeO.sub.2 film and aluminum is used for the electrodes has recently been reported. FIG. 1 is a sectional view showing a conventional MFIS FET using a CeO.sub.2 film.\nReferring to FIG. 1, the MFIS FET includes a p-type silicon substrate 1, a field oxide film 2, a source/drain region 3, a CeO.sub.2 film 5, a PbTiO.sub.3 ferroelectric film 6, and an aluminum gate electrode 7. In this device, the hetero-epitaxial growth of PbTiO.sub.3 ferroelectric film 6 is realized by using the CeO.sub.2 film 5 as a gate insulating film. Since the forming of the gate electrode in this structure occurs after high-temperature process of forming the ferroelectric film, aluminum can be used as the electrode material and the structure becomes suitable for integration. Thus, the MFIS FET can be applied to a memory device.\nHowever, although the CeO.sub.2 film 5 shown in FIG. 1 exhibits little lattice mismatching (approximately 0.35%) with a (111) substrate, it has no great advantage with respect to a (100) substrate. Furthermore, the electrical characteristics of this MFIS FET have not been completely verified."}
{"text": "The present invention relates generally to a connector and, more particularly, to a unique tool for removing contacts from a connector body.\nWhile the present invention will be described specifically in connection with a tool for removing contacts from an electrical connector, it will be appreciated that the tool could also be utilized for removing ferrules from fiber optic connectors.\nIn U.S. Pat. No. 3,110,093 to Johnson, there is disclosed an electrical connector of the type wherein one or more contacts, each connected to a wire, are inserted from the rear into a contact receiving bore in a connector body or insulator after the connector has been otherwise completely fabricated or assembled. The insulator includes a retaining clip or other locking means between the individual contacts and their respective bore walls for retaining the contacts in their operative position in the insulator. The Johnson patent also discloses a suitable tool for insertion into clearance space between the contact and the bore wall from the rear of the insulator to disengage the locking device and thereby permit manual withdrawal of the contact from the rear of the insulator by pulling on the wire attached to the contact. While this arrangement works well when a wire is present, the tool is ineffective fr an unwired contact in that the wire receiving portion of the contact normally does not extend beyond the rear of the insulator to permit the contact to be gripped for withdrawal.\nU.S. Pat. No. 3,380,141 Rofer discloses a tool for removing an unwired contact which, when inserted into the contact bore from the rear of the insulator, releases the locking means therein and frictionally engages the rear of the contact so that it may be withdrawn from the bore.\nFrench Pat. No. 2,240,600 discloses a front release contact retention assembly, in contrast to the rear release arrangments disclosed in the aforementioned Johnson and Rofer patents, wherein a contact extraction tool is inserted through the contact bore from the front of the insulator to release the contact locking means therein. The tip of the tool is longitudinally slotted and dimensioned so that it will frictionally engage the contact. Thus, as in the Rofer arrangement, the contact locking means is released while the tool frictionally engages the contact, thereby allowing the contact to be removed from the insulator. If manufacturing tolerances are not closely held in the forming of the prior art tools, it will be appreciated that adequate frictional engagement of the tool with the contact may not occur with the result that the tip of the tool may slip off the contact when the tip is withdrawn from the contact bore. Furthermore, even if the tool initially frictionally engages the contact with adequate force to allow the contact to be withdrawn from the insulator, continued use of the tool can result in the tip of the tool wearing so that it no longer firmly engages the contact.\nIt is the object of the present invention to provide a contact extraction tool which embodies manually operable gripping jaws which will assure that a firm gripping force can be applied to the contact so that the contact will be pulled out of the insulator when the tip of the tool is withdrawn from the contact bore."}
{"text": "A data center is a facility that physically houses various equipment, such as computers, servers (e.g., web servers, application servers, database servers), switches, routers, data storage devices, load balancers, wire cages or closets, vaults, racks, and related equipment for the purpose of storing, managing, processing, and exchanging data and information between hosts. A host typically corresponds to a computer system, a processor with associated physical memory, or a processor with shared memory and/or resources. Data centers also provide application services and management for various data processing functions, such as web hosting internet, intranet, telecommunication, and information technology.\nData centers are a unique environment because all the machines and services provided to clients are within a controlled and well-monitored environment. Specifically, data centers are managed by one or more administrators. In order to provide the application services and data processing functions, software is installed by an administrator on each host of the data center. Typically, installing the software on each host is performed remotely. Provisioning is the act of installing the software. One type of software that is typically provisioned is the operating system (OS) on bare hardware.\nProvisioning an OS in a data center is challenging because data centers are typically large and heterogeneous with respect to different operating systems (e.g., Solaris™ developed by Sun Microsystems, Inc. (a trademark of Sun Microsystems, Inc. located in Santa Clara, Calif.), Redhat Linux® developed by Red Hat, Inc. (a registered trademark of Red Hat, Inc. located in Raleigh, N.C.), SUSE® Linux developed by Novell, Inc. (a registered trademark of Novell, Inc. located in Waltham, Mass.), Windows® developed by Microsoft Corporation (a registered trademark of the Microsoft Corporation located in Redmond, Wash.), etc.) and types of hosts (i.e., different hosts use different hardware).\nFurther, each OS has a separate provisioning program that orchestrates the provisioning of the OS on the host. Each provisioning program typically has different input parameters, configuration file formats, and configuration steps. Accordingly, provisioning the different types of operating systems requires the administrator to be knowledgeable about the different input parameters, configuration file formats, and configuration steps."}
{"text": "This invention relates to the expression of fructose 1,6 bisphosphate aldolase (FDA) in transgenic plants to increase or improve plant growth and development, yield, vigor, stress tolerance, carbon allocation and storage into various storage pools, and distribution of starch. Transgenic plants expressing FDA have increased carbon assimilation, export and storage in plant source and sink organs, which results in growth, yield and quality improvements in crop plants.\nRecent advances in genetic engineering have provided the prerequisite tools to transform plants to contain alien (often referred to as xe2x80x9cheterologousxe2x80x9d) or improved endogenous genes. These genes can lead either to an improvement of an already existing pathway in plant tissues or to an introduction of a novel pathway to modify product levels, increase metabolic efficiency, and or save on energy cost to the cell. It is presently possible to produce plants with unique physiological and biochemical traits and characteristics of high agronomic and crop processing importance. Traits that play an essential role in plant growth and development, crop yield potential and stability, and crop quality and composition include enhanced carbon assimilation, efficient carbon storage, and increased carbon export and partitioning.\nAtmospheric carbon fixation (photosynthesis) by plants represents the major source of energy to support processes in all living organisms. The primary sites of photosynthetic activity, generally referred to as xe2x80x9csource organsxe2x80x9d, are mature leaves and, to a lesser extent, green stems. The major carbon products of source leaves are starch, which represents the transitory storage form of carbohydrate in the chloroplast, and sucrose, which represents the predominant form of carbon transport in higher plants. Other plant parts named xe2x80x9csink organsxe2x80x9d (e.g., roots, fruit, flowers, seeds, tubers, and bulbs) are generally not autotrophic and depend on import of sucrose or other major translocatable carbohydrates for their growth and development. The storage sinks deposit the imported metabolites as sucrose and other oligosaccharides, starch and other polysaccharides, proteins, and triglycerides.\nIn leaves, the primary products of the Calvin Cycle (the biochemical pathway leading to carbon assimilation) are glyceraldehyde 3-phosphate (G3P) and dihydroxyacetone phosphate (DHAP), also known as triose phosphates (triose-P). The condensation of G3P and DHAP into fructose 1,6 bisphosphate (FBP) is catalyzed reversibly by the enzyme fructose 1,6 bisphosphate aldolase (FDA), and various isozymes are known. The acidic isoenzyme appears to be chloroplastic and comprises about 85% of the total leaf aldolase activity. The basic isoenzyme is cytosolic. Both isoenzymes appear to be encoded by the nuclear genome and are encoded by different genes (Lebherz et al., 1984).\nIn the leaf, the chloroplast FDA is an essential enzyme in the Calvin Cycle, where its activity generates metabolites for starch biosynthesis. Removal of more than 40% of the plastidic aldolase enzymatic activity by antisense technology reduced leaf starch accumulation as well as soluble proteins and chlorophyll levels but also reduced plant growth and root formation (Sonnewald et al., 1994). In contrast, the cytosolic FDA is part of the sucrose biosynthetic pathway where it catalyzes the reaction of FBP production. Moreover, cytosolic FDA is also a key enzyme in the glycolytic and gluconeogenesis pathways in both source and sink plant tissues.\nIn the potato industry, production of higher starch and uniform solids tubers is highly desirable and valuable. The current potato varieties that are used for french fry production, such as Russet Burbank and Shepody, suffer from a non-uniform deposition of solids between the tuber pith (inner core) and the cortex (outer core). French fry strips that are taken from pith tissue are higher in water content when compared to outer cortex french fry strips; cortex tissue typically displays a solids level of twenty-four percent whereas pith tissue typically displays a solids level of seventeen percent. Consequently, in the french fry production process, the pith strips need to be blanched, dried, and par-fried for longer times to eliminate the excess water. Adequate processing of the pith fries results in the over-cooking of fries from the high solids cortex. The blanching, drying, and par frying times of the french fry processor need to be adjusted accordingly to accommodate the low solids pith strips and the high solids cortex strips. A higher solids potato with a more uniform distribution of starch from pith to cortex would allow for a more uniform finished fry product, with higher plant throughput and cost savings due to reduced blanch, dry and par-fry times.\nAlthough various fructose 1,6 bisphosphate aldolases have been previously characterized, it has been discovered that overexpression of the enzyme in a transgenic plant provides advantageous results in the plant such as increasing the assimilation, export and storage of carbon, increasing the production of oils and/or proteins in the plant and improving tuber solids uniformity.\nThe present invention provides structural DNA constructs that encode a fructose 1,6 bisphosphate aldolase (FDA) enzyme and that are useful in increasing carbon assimilation, export, and storage in plants.\nIn accomplishing the foregoing, there is provided, in accordance with one aspect of the present invention, a method of producing genetically transformed plants that have elevated carbon assimilation, storage, export, and improved solids uniformity comprising the steps of:\n(a) Inserting into the genome of a plant a recombinant, double-stranded DNA molecule comprising\n(i) a promoter that functions in the cells of a target plant tissue,\n(ii) a structural DNA sequence that causes the production of an RNA sequence that encodes a fructose 1,6 bisphosphate aldolase enzyme,\n(iii) a 3xe2x80x2 non-translated DNA sequence that functions in plant cells to cause transcriptional termination and the addition of polyadenylated nucleotides to the 3xe2x80x2 end of the RNA sequence;\n(b) obtaining transformed plant cells; and\n(c) regenerating from transformed plant cells genetically transformed plants that have elevated FDA activity.\nIn another aspect of the present invention there is provided a recombinant, double-stranded DNA molecule comprising in sequence\n(i) a promoter that functions in the cells of a target plant tissue,\n(ii) a structural DNA sequence that causes the production of an RNA sequence that encodes a fructose 1,6 bisphosphate aldolase enzyme,\n(iii) a 3xe2x80x2 non-translated DNA sequence that functions in plant cells to cause transcriptional termination and the addition of polyadenylated nucleotides to the 3xe2x80x2 end of the RNA sequence.\nIn a further aspect of the present invention, the structural DNA sequence that causes the production of an RNA sequence that encodes a fructose 1,6 bisphosphate aldolase enzyme is coupled with a chloroplast transit peptide to facilitate transport of the enzyme to the plastid.\nIn accordance with the present invention, an improved means for increasing carbon assimilation, storage and export in the source tissues of various plants is provided. Further means of improved carbon accumulation in sinks (such as roots, tubers, seeds, stems, and bulbs) are provided, thus increasing the size of various sinks (larger roots, tubers, etc.) and subsequently increasing yield and crop productivity. The increased carbon availability to these sinks would also improve composition and use efficiency in the sink (oil, protein, starch and/or sucrose production, and/or solids uniformity).\nVarious advantages may be achieved by the aims of the present invention, including:\nFirst, increasing the expression of the FDA enzyme in the chloroplast would increase the flow of carbon through the Calvin Cycle and increase atmospheric carbon assimilation during early photoperiod. This would result in an increase in photosynthetic efficiency and an increase in chloroplast starch production (a leaf carbon storage form degraded during periods when photosynthesis is low or absent). Both of these responses would lead to an increase in sucrose production by the leaf and a net increase in carbon export during a given photoperiod. This increase in source capacity is a desirable trait in crop plants and would lead to increased plant growth, storage ability, yield, vigor, and stress tolerance.\nSecond, increasing FDA expression in the cytosol of photosynthetic cells would lead to an increase in sucrose production and export out of source leaves. This increase in source capacity is a desirable trait in crop plants and would lead to increased plant growth, storage ability, yield, vigor, and stress tolerance.\nThird, expression of FDA in sink tissues can show several desirable traits, such as increased amino acid and/or fatty acid pools via increases in carbon flux through glycolysis (and thus pyruvate levels) in seeds or other sinks and increased starch levels as result of increased production of glucose 6-phosphate in seeds, roots, stems, and tubers where starch is a major storage nonstructural carbohydrate (reverse glycolysis). This increase in sink strength is a desirable trait in crop plants and would lead to increased plant growth, storage ability, yield, vigor, and stress tolerance.\nFourth, the invention is particularly desirable for use in the commercial production of foods derived from potatoes. Potatoes used for the production of french fries and other products suffer from a non-uniform distribution of solids between the tuber pith (inner core) and the cortex (outer core). Thus, french fry strips from the pith regions of such tubers have a low solids content and a high water content in comparison to cortex strips from the same tubers. Therefore, the french fry processor attempts to adjust the processing parameters so that the final inner strips are sufficiently cooked while the outer cortex strips are not overcooked. The results of such adjustments, however, are highly variable and may lead to poor quality product. Transgenic potatoes expressing fda will provide to the french-fry and potato chip processor a raw product that consistently displays a higher tuber solids uniformity with acceptable agronomic traits. In the french fry plant production process, inner pith fry strips from higher solids uniformity tubers will require less time to blanch, less time to dry to a specific solids content, and less time to par-fry before freezing and shipping to retail and institutional end-users.\nTherefore, with respect to potatoes, the present invention provides 1) a higher quality, more uniform finish fry product in which french fries from all tuber regions, when processed, are nearly the same, 2) a higher through-put in the french fry processing plant due to lower processing times, and 3) processor cost savings due to lower energy input required for lower blanch, dry, and par-fry times. A raw tuber product that displays a higher solids uniformity will also produce a potato chip that has a reduced saddle curl, and a reduced tendency for center bubble, which are undesirable qualities in the potato chip industry. Reduced fat content would also result; this would contribute to improved consumer appeal and lower oil use (and costs) for the processor. The increase in solids uniformity will also translate to an increase in overall tuber solids. For both the french fry and chipping industries, this overall tuber solids increase will also result in higher through-put in the processing plant due to lower processing times, and cost savings due to lower energy input for blanching, drying, par-frying, and finish frying."}
{"text": "Despite the fact that tap water having been filtered, settled, sterilized etc. and meeting the government standards for drinking water, the tap water in daily use is inevitably going through second contamination when it is carried in the pipeline due to the rust on the pipeline or external infiltration, which may bring excessive bacteria, viruses, residual chlorine, organic materials, heavy metal, radioactive substances and other various hazardous substances to the tap water. A water purification device is therefore a device that could filter out all these undesirable substances.\nSince access to hygienic drinking water has drawn public attention in recent years, water purification device is widely used in daily life. Various water purification devices purify water depending on filter cartridges; therefore a filter cartridge is the core of water purification devices. There is a need for replacement of the filter cartridge on a regular basis to keep better purifying effect. A high demand for water purification devices keeps increasing in the current market, however a large amount of counterfeiting filter cartridges arise. If the material quality of filter cartridges cannot be guaranteed, the purifying effect will be in adverse as opposed to using an original filter cartridge. The use of such counterfeiting filter cartridge cannot bring desirable purifying effect or meet standards for drinking water; instead, the operational life span of such counterfeiting filter cartridge is short, use of which is adverse to the health of consumers and the reputation of manufacturers."}
{"text": "The present invention relates generally to devices and methods for measuring haze and transmissivity in transparencies, and more particularly to a device and method for measuring haze and transmissivity of aircraft transparencies as viewed through night vision goggles.\nNight vision goggles (NVGs) are being used by aircrew members with increasing frequency during nighttime flight and ground operations. Concurrently, aircraft transparencies are constantly undergoing improvements in materials such as acrylic to polycarbonate and their surfaces are being treated with gold, indium-tin-oxide or other coatings. Ground vehicles are being equipped with bullet-resistant glass which may be fabricated using a combination of specially treated plastic and glass materials. These material changes may attenuate the infrared (IR) energy utilized by NVGs thereby degrading the visual acuity of the aircrew member. Aircrew members have informally reported lowered visual performance while using NVGs in some aircraft equipped with coated canopy systems. Traditional transmissivity and haze measurements made in the visible spectral region (400-700 nm) cannot characterize the entire problem. Preliminary laboratory measurements of sample coated canopies found that IR transmissivities of windscreen coatings varied considerably among manufacturers thereby affecting observer performance (see, Pinkus et al, \"The Effects of Aircraft Transparencies on Night Vision Goggle-Mediated Visual Acuity,\" SAFE Symposium Proceedings (1997) pp 93-104). However, NVG-mediated visual performance through canopies is not solely affected by IR attenuation, but also by haze resulting from scatter of incident light.\nThe invention solves or substantially reduces in critical importance problems in the prior art by providing a portable device and method for measuring transmissivity and haze of aircraft transparencies as viewed through NVGs. The invention includes a pair of light sources for projecting a transmitted beam and a haze producing beam onto a transparency and sensor portion including a light intensifier tube of an NVG and photometer for measuring the luminance output thereof in order to quantify the NVG weighted attenuation (transmissivity) and haze (light scatter) characteristics of the transparency.\nThe invention finds utility for measuring the IR transmissivity and IR haze characteristics of transparencies important in selection of coating materials (gold, indium-tin-oxide) for use with NVGs, quality control of the coatings within and among transparency manufacturers, quick-response field evaluations of transparency and NVG integration issues, life-cycle costs of coatings, monitoring deleterious environmental effects (acid rain, sand abrasion, bio-chemical) on coating integrity, maintenance procedures, and visual performance of an NVG-equipped aircrew member.\nIt is therefore a principal object of the invention to provide system and method for measuring transmissivity and haze in a transparency.\nIt is a further object of the invention to provide system and method for measuring infrared transmissivity and haze in a transparency as viewed using night vision goggles.\nIt is another object of the invention to provide system and method for measuring the effect of using night vision goggles on the observed infrared transmissivity and haze in a transparency.\nIt is another object of the invention to provide a device and method for field evaluation of transparencies.\nIt is yet another object of the invention to provide an improved system and method for measuring transmissivity and haze in aircraft transparencies.\nIt is yet another object of the invention to provide a reliable portable transmissivity and haze measurement device.\nThese and other objects of the invention will become apparent as a detailed description of representative embodiments proceeds."}
{"text": "1. Field of the Invention\nThe invention relates to medical patient weighing, and, more particularly, to a portable device, system and method for weighing ambulatory and non-ambulatory patients.\n2. Description of the Background\nIt is known that persons confined to a wheelchair often need to be weighed as part of a medical check up or follow-up. For many of these individuals, stepping off a wheelchair and onto a weight scale is difficult, or impossible, to accomplish. Consequently, weighing both the wheelchair and the patient, and subtracting the weight of the wheel chair, is a practical method of weighing. This practice has become commonplace. Such a simple procedure may, in fact, be quite involved, because the weighing device must accommodate the wheelchair, which generally requires ramps and a platform scale, for example. As such, wheelchair users ordinarily are weighed at a physician's office, hospital, nursing home, clinic, or the like, using a platform scale, which weighing, due to the inconvenience of the need to travel, may occur once a month or less. Such infrequent weighing renders weight control programs quite difficult to maintain. In addition to those using wheelchairs, often those requiring crutches or walkers must be seated in wheelchairs for weighing due to the difficulty in maintaining balance on an ordinary scale while utilizing crutches or a walker.\nDue to the size of platform scales, platform scales are generally found principally in clinics, hospitals, nursing homes, or rehabilitation centers with a significant number of mobility-disabled patients. Consequently, physicians rarely have office platform scales, not only because such scales occupy a disproportionate amount of floor space in relation to usage, but also because, even with foldable ramps, such scales occupy valuable storage space. Although physician's weighing chairs may be available with built-in scales, those confined to wheelchairs may find it inconvenient or dangerous to be transferred between chairs for a procedure as routine as weighing.\nIn response to the deficiencies inherent in available weighing systems for the mobility disabled, portable weighing devices including two laterally displaced, electrically connected and rigidly aligned weighing modules have been developed. An example of one such system is documented in U.S. Pat. No. 5,414,225. The weighing device taught therein is designed for the conventional manually-operated wheelchair supported by two front caster wheels and two rear main wheels. A notable feature of the above-referenced device is the rigidly adjustable scissors-like adjustment between the two weighing pads. The rigidly adjustable pads must be pre-extended to the correct width of the wheelchair before use. A continuously variable width adjustment would be attractive because it permits easy last-minute adjustment of the weighing pads by an operator.\nAdditionally, it would be advantageous if a weighing scale could be utilized by both non-ambulatory and ambulatory patients, wherein ambulatory patients could step on the scales and be weighed without the use of the wheelchair, and wherein non-ambulatory patients may be pushed onto the same scale in a wheelchair. Thus, dual purpose weighing pads could be used for both purposes, and a smaller platform scale in a doctor's office setting would not be needed.\nTherefore, a need exists for a portable, fully adjustable wheelchair scale that can additionally serve as a substitute for a platform scale."}
{"text": "1. Field of the Invention\nThis invention relates to intravascular implants for maintaining vascular patency in human blood vessels. More particularly, this invention relates to a radially expandable stent made from a fine wire formed into a serpentine ribbon wound into a cylindrical shape for introduction into a body vessel for balloon expansion therein in a radial fashion to support the wall of the vessel when in the expanded configuration. The stent includes an impermeable membrane that lies in the plane of the cylinder. This invention is particularly useful in transluminal implantation of a stent for use in the coronary angioplasty to prevent restenosis, and for the treatment of aneurysms or subintimal dissections.\n2. Description of the Prior Art\nThe basic concept of stents has been known for a number of years. Various types of stents have been proposed and patented, including self-expanding spring types, compressed spring types, mechanically actuated expandable devices, heat actuated expandable devices, and the like. More recently, expandable sleeves have been proposed such as shown in Palmaz, U.S. Pat. No. 4,733,665. In this disclosure there is shown a sleeve having slots therethrough to form a permeable mesh. The sleeve is placed transluminally, and then expanded by a balloon catheter through the elastic limit of the metal so as to permanently deform the sleeve into supporting contact with the interior surface of a blood vessel. Other examples of expandable wire stents are shown in Hillstead, U.S. Pat. No. 4,865,516, and Wiktor, U.S. Pat. No. 4,886,062.\nIn my U.S. Pat. No. 5,217,483 (referred to hereinafter as the '483 Patent) I disclose the use of a fine platinum wire bent into a serpentine flat ribbon which is wound around a mandrel into a radially expandable cylindrical sleeve for mounting on a balloon catheter for transluminal intravascular placement. The expanded sleeve as disposed in a vessel possesses gaps or interstices which are believed to protect the vascular endothelium from contact therewith, and to promote endothelial proliferation and remodeling of the vascular intima about the sleeve. However in many cases the trauma of implantation causes microscopic or even macroscopic intimal tears, and exposes the subintimal space to the bloodstream. This is undesirable as further dissection of the vessel is possible. Furthermore tissue reactions, including thrombus formation, intimal fibroplasia, and cicatrization could result in restenosis of the vessel.\nPermeable transluminally placed intravascular stents are suboptimum for the treatment of aneurysms. While the structure of the expanded mesh may reduce the pulsatile forces acting against a weakened arterial wall, this effect is doubtless incomplete. The art has long utilized grossly impermeable continuous structures such as Teflon grafts and implants for aneurysm repair in larger vessels by conventional open surgical techniques."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication system, especially to a method for transmitting SRS in an LTE communication system and an apparatus using the same.\n2. Description of the Related Art\n3GPP (The 3rd Generation Partner Project) standardization organization is working out a next generation of wireless communication standard which is named LTE (Long Term Evolution). In a physical layer interface, the new standard adopts OFDM (Orthogonal Frequency Division Multiplexing) technology, which is different from conventional CDMA (Code Division Multiple Access) technology. OFDMA is used in downlink and SCFDMA (Single Carrier Frequency Division Multiple Access) is used in uplink. The technology used in the new standard is effective to resist multi-path propagation, with the adoption of frequency domain equalization reducing complexity of the conventional time domain equalization, and is more suitable to bandwidth high-speed data transmission.\nFrom a point of view of air-interface, the LTE standard techniques can be divided into two categories: a TDD (Time Division Duplex) system and an FDD (Frequency Division Duplex) system. The LTE system supports variable bandwidths. And typical bandwidths include 1.4 MHz, 3 MHz, 5 MHz, 10 MHz, 15 MHz and 20 MHz, which can meet demands of different scenarios.\nFIG. 1 illustrates a physical layer frame structure for a LTE FDD system in which a length of radio frame (101) is 10 ms, consisting of ten equally sized radio sub-frames (102) of 1 ms length. Each radio sub-frame consists of two equally sized timeslots (103) of 0.5 ms length.\nFIG. 2 illustrates a PHY layer frame structure for LTE TDD system. As shown in FIG. 2, a length of radio frame (201) is 10 ms, consisting of ten equally sized radio sub-frames (204) of length 1 ms. Each five continuous radio sub-frames consists a half-frame (202) of length 5 ms. Different from the LTE-FDD system, a second (211) and seventh (212) radio sub-frame in LTE-TDD radio frame are two special sub-frames. A length of the special sub-frame is 1 ms, consisting of three special slots, indicating DwPTS (205 or 208), GP (206 or 209) and UpPTS (207 or 210) respectively. The lengths of the three special slots are variable and are defined by system, and the total length is 1 ms. The length of UpPTS can be 0, 1 or 2 SCFDMA symbols. If the length of UpPTS is 2, UpPTS is used to transmit the uplink Short RACH or Uplink SRS signal or both the Short RACH and SRS signal. If the length of UpPTS is 1, UpPTS is used to transmit the uplink SRS signal. The other eight sub-frames except the special two are respectively consist of two slots (203) of length 0.5 ms.\nIn the LTE system, according to network scheduling, UE (User Equipment) sends an SRS (Sounding Reference Signal) to eNodeB (evolved NodeB). The SRS signal is used to: according to an analysis result of the SRS signal, eNodeB estimates a quality of channel which is used for transmitting SRS from UE to eNodeB and scheduling data according to frequency selective characteristics; eNodeB performs timing tracking for UE by analyzing the SRS signal and performs a close-loop power control. According to a current standardizing process, main conclusions for SRS transmission in LTE FDD system include: eNodeB broadcasts the SRS in a designated cell as needed and SRS is transmitted in some sub-frame in a designated cell periodically. A period is selected from {2, 5, 10, 20, 40, 80, 160, 320} ms. After the UE receives the SRS in the designated cell, the ODFM symbol resource occupied by the SRS is not used when transmitting uplink data. In order to perform the transmission of SRS, UE should receive a user-designated SRS signal transmitted from the network. The signal informs the user of the OFDM symbol resource used to transmit SRS. Currently, there is no description for transmitting the SRS of designated UE in a PHY layer specification which is accomplished in LTE.\nNowadays, a basic idea in the standard for the user-designated SRS signaling is that the signaling includes three parts: Duration, Period and Offset in which, the duration can use 1 bit to indicate that just one snapshot or infinite. The period value is selected from {2, 5, 10, 20, 40, 80, 160, 320} ms. In LTE FDD, the offset is a time between twice transmission time of each OFDM symbol of SRS from the beginning of the SRS period, and the basic unit is 1 ms. In LTE TDD, definition of offset is different from that in LTE FDD. Since in LTE TDD, SRS can be transmitted in UpPTS or the other uplink sub-frame, the uplink sub-frame may be discontinuous and UpPTS occupy two OFDM symbols at most, the offset is defined as an interval between an OFDM symbol position used to transmit SRS and an OFDM symbol position used to transmit SRS until the period of SRS transmission starts. For example, if the SRS symbol position at the period beginning is defined as 0, a symbol position used to transmit SRS is 3 means that the interval between the two symbols is 3, there is at most 2 OFDM symbol positions can be used to transmit SRS.\nThe manner of transmitting SRS in LTE TDD is mainly the same as that in LTE FDD. However the system structure of LTE TDD is different from of that in LTE FDD. Difference is that in LTE TDD, a half-frame of length 5 ms has both uplink sub-frame and downlink sub-frame, a number of uplink sub-frames and downlink sub-frames is configured by the network. In some configuration, a half-frame of length 5 ms at least has one uplink sub-frame (exclude UpPTS). According to a principle that only one SRS is transmitted in one uplink sub-frame, there is only one SRS transmission in every 5 ms, and the system can't achieve the SRS transmission with a 2 ms period. Therefore, the performance of SRS transmission by UE is deteriorated in a fast Time-varying channel.\nBased on the difference between the LTE TDD and LTE FDD, the current configuration of the 2 ms transmission period for SRS in the LTE FDD cannot be used in the LTE TDD system."}
{"text": "The wavelength of light for lithography has been reduced into the deep ultraviolet (DUV) range to produce the feature size necessary for current and future electronics devices. The electronics industry is developing new resists that are tailored to the DUV range. One such resist class is chemically amplified resists.\nThe main components of chemically amplified resist formulations are a photoacid generator compound, a polymer resin and a solvent capable of dissolving the photogenerator and the resin. For many positive chemically amplified resists, the polymer resin contains acid labile groups which makes the polymer resin insoluble in an aqueous developer. Upon irradiation, the photoacid generator compound produces an acid which cleaves the acid labile groups resulting in a polymer resin that is aqueous soluble. Chemically amplified resists have generated a great deal of interest and there are numerous patents available discussing these compositions such as, for example, U.S. Pat. Nos. 5,069,997; 5,035,979; 5,670,299; 5,558,978; 5,468,589; and 5,389,494.\nOne group of polymers which can be used as resins in chemically amplified resists are acetal derivatized polymers. The alkali solubility of phenolic resins are greatly inhibited by converting the hydroxyl groups to acetal groups. Typically, acetal phenolic resins are produced by reacting a phenolic resin with a vinyl ether in the presence of an acid catalyst.\nThe acetal resins are then formulated with a photoacid generator compound and solvent to form a chemically amplified resist product. Upon irradiation, the generated acid cleaves the acetal groups, and a phenolic resin is created which is soluble in an aqueous developer.\nOne problem with producing acetal derivatized polymers is that currently there are only a very limited number of bulk vinyl ethers that can be used to produce the acetals. It is possible to generate other acetal functionalities by synthesizing intermediates, but this would require a series of reactions which makes the overall process for producing acetal polymers relatively complicated, expensive, and unlikely to be reproducible. Thus, the limited number of bulk vinyl ethers severely constrains the use of acetal polymers because only a small number of acetal groups can be readily and inexpensively generated.\nIt would be advantageous to be able to generate reproducibly and economically a large variety of acetal groups on polymers. Various resist properties such as alkaline solubility, etch resistance, film shrinkage, temperature stability, adhesion, sensitivity can be altered by choosing appropriate acetal functionalities. Furthermore, if more than one acetal group can be readily generated on the polymer, this would provide further ability to tailor the polymer resin to a specific application. For example, one acetal group may be used to increase the sensitivity, while another acetal group may be used to increase the etch resistance. The host of applications for acetal polymers can be greatly expanded if a large variety of acetal groups on the polymers can be produced readily.\nThe present inventors have developed a reproducible process for readily producing a wide variety of acetal derivatized polymers.\nIt is another object of this invention to provide a new, inexpensive, reproducible process for generating a large variety of acetal derivatized polymers that have applications as resins for photoresists.\nIt is a further object of this invention to provide a process for generating a large variety of acetal derivatized polymers based on mixed acetals. A still further object of this invention is to provide a variety of mixed acetals derivatized polymers that have applications as resins for photoresist. Yet another object of this invention is to provide new photoresist compositions containing mixed acetal derivatized polymers and the use of such new photoresist compositions in photolithography imaging processes to produce microelectronic devices. The present invention has many important advantages over the prior art. First, the present invention provides a new class of polymers based on mixed acetals. Desired properties of the polymer resin can be tailored by choosing among appropriate acetal functionalities. Also, the relative proportions of the different acetals in the polymer can be readily varied by changing the relative proportions of reagents in the feedstock. Changing the proportion of the mixed acetals can further tailor the properties of the polymer resin. In addition, the reproducibility of the mixed acetal polymers compositions are excellent. Reproducible polymer compositions are very important in order to produce commercially viable resist formulations because the properties of the resist must not change from batch to batch. Furthermore, the mixed acetal process is relatively inexpensive because intermediates needed for the reaction are commercially available and only one synthesis reaction is required to produce the final acetal polymer. Previously, a number of synthesis reactions would generally be required to produce the final single acetal polymer if the acetal group was not from one of the readily available vinyl ethers.\nThe present invention also provides many additional advantages which shall become apparent as described below."}
{"text": "This invention relates to engine and vehicle control through feedback of driving performance.\nAutomotive vehicles are ultimately controlled by a human operator that makes decisions about the vehicle including speed, direction, and operating parameters such as transmission gear, accessory operation, refueling, maintenance, and the like. While many operators develop a general feel for how a vehicle is performing, the vehicle generally provides little if any quantitative information to the driver indicating how efficiently the vehicle is operating.\nAn example of the need for vehicle efficiency feedback occurs with on-highway heavy trucks. Vehicle operators have direct or indirect control over a wide range of vehicle parameters including engine fuel, service (wheel) brakes, engine retarders such as engine brakes and Jake brakes, transmission retarders, transmission gear, idling time, use of engine driven accessories such as air conditioning, and the like. Each of these parameters affects the efficiency of operation and, hence, the cost of operation. Truck owner/operators and fleet owners are highly concerned with operational costs and are therefore in need of quantitative measures of vehicle operating efficiency.\nOne problem with previous measures of vehicle operating efficiency is a failure to provide meaningful cost-related information relevant to vehicle control. Measures, such as instantaneous miles per gallon of fuel consumed (MPG), tend to relay inaccurate cost indicators to a vehicle driver or automatic controller. For example, MPG decreases during vehicle acceleration and increases during vehicle braking. This often does not reflect optimum driving requirements. Acceleration adds energy to the vehicle and braking removes energy from the vehicle. Braking, in particular, should be avoided if possible to prevent loss of vehicle energy through heat.\nIt is an object of the present invention to provide quantitative measures of vehicle operating efficiency. These measures may be used for a variety of purposes including display to the vehicle operator, storage in memory for future analysis, and input to engine controllers.\nIn carrying out the objects and features of the present invention, a method is provided for controlling an automotive vehicle. At least one vehicle parameter is sampled. At least one measure of driving efficiency is determined based on the sampled vehicle parameter. An estimate of operating cost is found based on the at least one measure of driving efficiency.\nIn an embodiment of the present invention, at least one driving efficiency limit is set. An increase in the maximum allowed vehicle velocity is permitted if the measure of driving efficiency exceeds the driving efficiency limit.\nIn embodiments of the present invention, measures of driving efficiency may include measures of normalized cost, measures of shifting and idling cost, and measures of change in vehicle kinetic energy such as, for example, may be caused by vehicle braking.\nA method for gauging vehicle efficiency is also provided by the present invention. Vehicle velocity is repetitively sampled. If the sampled vehicle velocity is less than a previously sampled vehicle velocity and if braking is applied, the sampled braking velocity cost is determined based on the difference between the square of the sampled vehicle velocity and the square of the previously sampled vehicle velocity. A measure of driving efficiency is then determined based on the sampled braking velocity cost. These measures may include instantaneous braking velocity cost per mile and average braking velocity cost per mile.\nIn another method for gauging driving efficiency provided by the present invention, an engine revolution idle cost is determined. The engine speed is repetitively sampled. At least one sampled shifting and idling cost is determined as the product of a sampled engine speed and the engine revolution idle cost. A measure of driving efficiency is found based on the sampled shifting and idling cost. These measures may include instantaneous shifting and idling cost per mile and average shifting and idling cost per mile.\nA driving efficiency gauge is also provided. A display presents at least one measure of driving efficiency to the vehicle operator. A memory holds a sequence of at least one measure of driving efficiency. Control logic in communication with the electronic control module controlling an internal combustion engine receives samples of at least one vehicle operating parameter from the electronic control module. The sequence of driving efficiency measures is determined based on the vehicle operating parameter samples. At least one of the efficiency measures is displayed and the sequence of efficiency measures is stored in the memory.\nThe above objects and other objects, features, and advantages of the present invention will be readily appreciated by one of ordinary skill in the art from the following detailed description of the best mode for carrying out the invention when taken in connection with the accompanying drawings."}
{"text": "None\n1. Background-Field of Invention\nThe novel invention is a sliding stabilizer that operates in tandem with a mining roof support machine and is particularly useful in longwall mining operations.\n2. Background-Description of the Prior Art\nAs mineral mining becomes deeper, the threat of rock bursts becomes greater. This is a problem that has been common in deep hard rock and mineral mines for years. It will be an ever-increasing problem as mining at greater depths is required. Despite gains, this threat is not controlled. In deeper seams, consideration must be given to support, safety anti stability. Additional design is required. Currently three types of mining roof support machines are popular: frame, chock, and shield. All three have to a greater or lesser degree safety, stability and ventilation problems. Problems of particular consideration are: (1) the rock burst problem that exists from the unsupported roof created when a shearer passes a support machine; (2) the rock burst problem that exists as a roof covering of a support machine is lowered to advance; (3) lack of stability caused by oversize roof coverings in an attempt to create more ventilation; (4) mineral bursts from a mineral face; (5) instability in support machines caused by inclined seams; and (6) lack of work space in thin seams.\nRigid canopies have been designed in order to help with the above listed problems. They have not worked. (1) When the roof covering of the support machine is lowered, the rigid canopy also lowers., This creates a large unsupported roof and the potential for rock bursts. (2) The rigid canopy extends from the support machine and becomes unstable. Rock bursts can cause canopy rotation.\nIn three prior designs rigid canopies also have attached legs. These have not worked. (1) The legs interfere with the working space and ventilation. (2) The connections of the canopies break if the leg is not lowered with the roof covering of the support machine.\n3. Objects and Advantages\nMiners throughout the world are demanding safe working conditions. World energy demand continues to increase. Therefore, it is a general object and advantage of a sliding stabilizer to provide support, stability, safety, and efficiency in longwall mining operations.\nIn addition to the object and advantage above, it its the object and advantage of a sliding stabilizer:\n(1) to provide ventilation and working space; PA1 (2) to provide support, stability and safety to enable mining of inclined seams; PA1 (3) to provide ventilation and working space to enable mining of thin seams; PA1 (4) to provide support, stability and safety to enable mining of deeper seams; and PA1 (5) to provide support, stability, safety, ventilation and working space in the working zone of the shearer in longwall mining operations. PA1 (1) to prevent rock bursts from the unsupported roof created as a shearer passes a support machine; PA1 (2) to prevent rock bursts from the unsupported roof created when the roof covering of a support machine is lowered to allow the support machine to advance; PA1 (3) to prevent mineral bursts from the mineral face; PA1 (4) to prevent canopy rotation from rock bursts; and PA1 (5) to prevent canopy breakage when the roof covering of the support machine is lowered.\nIn addition it is the object and advantage of a sliding stabilizer:"}
{"text": "Atmospheric pollution is a societal problem which is receiving much attention. The major source of such pollution is the extensive use of fossil fuels, although industrial and chemical processes, such as the manufacture of nitric acid, also contribute. The principal pollutants are nitrogen oxides, carbon monoxide, and perhaps to a lesser extent hydrocarbons, sulfur oxides and other objectionable gases and vapors.\nAlthough several nitrogen oxides are known which are relatively stable at ambient conditions, it is generally recognized that two of these, viz. nitric oxide (NO) and nitrogen dioxide (NO.sub.2), are the principle contributors to smog and other undesirable environmental effects when they are discharged into the atmosphere. These effects will not be discussed further herein since they are well recognized and have led various government authorities to restrict industrial and automotive emissions in an attempt to limit the level of the oxides in the atmosphere. Nitric oxide and nitrogen dioxide, under appropriate conditions, are interconvertible according to the equation: EQU 2NO+O.sub.2 =2NO.sub.2.\nFor purposes of the present invention, NO.sub.x will be used herein to represent nitric oxide, nitrogen dioxide, and mixtures thereof\nFormation of nitrogen oxides from the elements occurs in the high temperature zones of combustion processes. The internal combustion engine, and coal or gas-fired furnaces, boilers and incinerators, all contribute to NO.sub.x emissions. In general, fuel-rich combustion mixtures produce exhaust gases with lower contents of NO.sub.x than do lean mixtures. Although the concentrations of NO.sub.x in the exhaust gases produced by combustion usually are low, the aggregate amounts discharged in industrial and/or highly populated areas is adequate to cause problems. Other industrial sources of pollution also exist. These are associated with the manufacture of nitric acid, with nitration of organic chemicals, and with other chemical operations, such as the reprocessing of spent nuclear fuel rods by dissolution in nitric acid to recover uranyl nitrate followed by calcination to convert the nitrate to uranium oxide. In these instances the exhaust gas may contain relatively high levels of NO.sub.x, such as from 0.1% up to 2% to 3%.\nThe so-called \"stable\" nitrogen oxides have in common the somewhat peculiar property that although they are thermodynamically extremely unstable at room temperature with respect to decomposition into elemental oxygen and nitrogen, no simple, economical method has been described for inducing this decomposition. It has been discovered, however, that adding a reductant such as ammonia to the exhaust gas can, under appropriate reaction conditions, convert NO.sub.x to elemental nitrogen and steam at moderate temperatures in the range of 200.degree. 600.degree.600.degree. C.\nUnlike the objectionable nitrogen oxides, carbon monoxide is a combustible gas which is convertible to CO.sub.2 by burning. Also unlike the nitrogen oxides, for which relatively few effective conversion catalysts are known, the combustion of carbon monoxide is catalyzed by a fairly large number of catalysts among which are the platinum group metals and their alloys, gold, base metal oxides such as chromia and copper oxide, and certain rare earth oxides.\nThe technology for abatement of pollution by noxious gases began to be developed some years ago. The earliest methods for control of gaseous emissions included absorption, adsorption, condensation, chemical reaction and incineration. Some of these methods, such as absorption, require the disposal of recovered pollutants, which in itself may be a problem.\nIn recent years, increasing attention has been paid to developing catalytic technology for control of exhaust gas pollution. Such technology potentially offer the advantages of low cost and the possibility for converting the noxious gas to one or more innocuous substances such as carbon dioxide and water, which are free of disposal problems. In instances in which two or more substances such as NO.sub.x and CO are present in the exhaust, as is often the case with fossil fuel combustion, a single catalyst that is effective for converting more than one of the pollutants is highly desirable.\nCertain zeolitic catalysts have been described as effective for control of pollution by one or more of NO.sub.x, CO, and hydrocarbons. The term \"zeolite\" or \"zeolitic\" as used herein refers to porous crystalline minerals or synthetic oxides, usually aluminosilicates, that have a rigid three dimensional framework structure, such as are described in \"Zeolite Molecular Sieves\" by Donald W. Breck, John Wiley & Sons, New York, N.Y. (1974), relevant portions of which are incorporated herein by reference for background.\nU.S. Pat. No. 3,900,554 to Lyon describes a non-catalyzed homogeneous gas phase reaction to remove NO.sub.x from combustion effluent by adding 0.4 to 10 moles (preferably 0.5 to 1.5 moles) of ammonia followed by heating to 1600.degree. C. to 2000.degree. C. The NO.sub.x content is lowered as a result of its being reduced to nitrogen by reaction with ammonia. The method is reported to work best if hydrocarbon is also added to the mixture. The extremely high temperature required for NO.sub.x abatement is a disadvantage of the method.\nU.S. Pat. No. 4,220,632 to Pence et al. discloses a process for reducing noxious nitrogen oxides from a fossil-fuel-fired power generation plant, or from other industrial plant off-gas stream, to elemental nitrogen and/or innocuous nitrogen oxides employing ammonia as reductant and, as catalyst, the hydrogen or sodium form of a zeolite having pore openings of about 3 to 10 Angstroms. The Pence et al. patent illustrates the so-called Selective Catalytic Reduction Process (hereinafter referred to as the \"SCR Process\") for removal of NO.sub.x from exhaust gas. While the SCR process usually operates at 200.degree.-600.degree. C., which is an advantage over the uncatalyzed conversion, an external reductant (ammonia) must be furnished which adds to the cost of the process, and control of the NH.sub.3 feed is required to avoid ammonia emissions.\nDuring early commercial development of molecular sieves, it was found that copper, for example, could be introduced into the inner absorption regions of the crystalline zeolites by several different means, including simple ion-exchange with divalent copper ion, and impregnation with complexes such as copper acetylacetonate in which the metal is in the zero valence state. It was further found that copper, in the inner absorption regions of the zeolite, whether present as metal or cation, reacted readily with molecules that could enter the pores of the molecular sieve. U.S. Pat. No. 3,013,985 to Breck et al. is incorporated herein by reference for background purposes, and also for the description contained therein on methods for loading the molecular sieve with copper either by ion-exchange or by the use of complex decomposable compounds in which the metal is in the zero valence state.\nThe use of copper exchanged zeolites in the SCR reduction of NO.sub.x is mentioned in U.S. Pat. No. 4,046,088.\nU.S. Pat. No. 3,346,328 to Sergeys et al. proposes to use as the catalyst for treating internal combustion engine exhaust gases a Cu.sup.++ exchanged zeolite such as Zeolite Y which has CuO loaded or held interstitially in its pore system. Since the copper ion is reported to play an important part in hydrocarbon conversion and copper oxide is reported to be an influential factor in carbon monoxide conversion, the activity of the catalyst is optimized for a particular exhaust gas by varying the proportion of zeolitic Cu.sup.++ cations and impregnated CuO.\nIn catalytic pollution control technology, the monetary cost is significant, and particularly so when the exhaust contains little or no economically recoverable values in terms of materials or heat. For this reason, it becomes very important to find catalysts that are relatively inexpensive, and that are not only effective for the desired conversion, but that also are adequately long lived in the process environment. As a general rule, a catalyst will deteriorate with time on stream for one or more reasons, such as loss of activity for conversion, for selectivity, or for both, until it becomes necessary to discard it. The term \"aging\", as used herein, will refer to such deterioration. One of the factors often responsible for excessive aging is exposure of the catalyst to steam at high temperature. Unfortunately, the burning of fossil fuels inherently forms a flue gas or exhaust gas that contains a significant content of steam as well as NO.sub.x and/or other pollutants such as CO and hydrocarbons. Catalytic removal of such pollutants necessarily include contact with steam.\nIt is an object of this invention to provide a steam-resistant zeolite copper catalyst for the catalytic decomposition of NO.sub.x. It is a further object of this invention to provide a method for preparing a steam-resistant zeolitic copper catalyst. These and other objects will become evident on reading this entire specification including the appended claims.\nWe now have found a surprising, simple way to make zeolitic copper catalysts that are highly resistant to aging in the presence of steam."}
{"text": "In recent years, various network-type image display systems have been developed with popularization of network technologies. For example, in a network-type projector system, images can be projected by transmitting image data to an arbitrary projector from another arbitrary device connected to a network, for example, from a terminal.\nIn the network-type image display system, the setup of network information and connecting operations of devices are complicated, which were difficult operations particularly for entry-level users. Accordingly, as a method of reducing difficulty of these operations, for example, a technique is disclosed (refer to JP-A-2005-184314), in which a wireless LAN card for the terminal is inserted into a slot of a wireless access point to thereby register a MAC address of the wireless LAN card in the access point information such as an SSID is written in the wireless LAN card, and the wireless LAN card is inserted into a slot of a terminal to thereby perform network setup between the access point and the terminal easily. As another example, a configuration is disclosed (refer to JP-2006-227791), in which a communication device having communication functions is attached to a data projector to thereby read setup of communication of the data projector, and the communication device is re-connected to a computer to enable transmission of images from the computer to the data projector.\nHowever, in the above techniques, devices having communication functions such as the wireless LAN card and the communication device are used, therefore, general versatility is low. Additionally, a problem that functions overlap when the computer itself has communication functions has been pointed out."}
{"text": "The present invention relates to composite plastic moldings having a solid, thin skin and a foamed backing layer and a process for their production. Such moldings are widely used in the manufacture of interior parts for motor vehicles, such as instrument panels, center consoles, arm rests or door inner panels.\nIn the production of such composite plastic moldings, the solid skin is generally produced in a known manner from a thermoplast in a first step. This step is carried out in a suitable mold, for example by deep drawing of an ABS film or by slush molding from a PVC plasticizer plastisol or a PVC powder containing polymeric plasticizers. A PVC skin is often preferred to an ABS skin by virtue of its more pleasant feel.\nIn a second step, the solid skin produced in the first step is back-foamed with a suitable plastic either in the same mold (deep drawing) or after transfer to a second mold (slush molding). By virtue of their excellent foamability and the range of variation of the physical properties of the foam, polyurethane(urea) systems are generally used for this back-foaming step.\nUnder certain loads, however, composite plastic moldings of this type are affected by serious problems, particularly where the skin is based on PVC. The PVC contains stabilizers and plasticizers which have a certain vapor pressure and, hence, can lead to the emission of unpleasant odors and to a permanent coating on the windows of the automobile from inside.\nIn addition, the effect of heat on the composite material (dashboards can become heated to surface temperatures of &gt;100.degree. C., depending upon color and insulation) produces considerable interaction between the outer skin and the foam backing, which can lead to delamination, embrittlement of the outer skin and degradation of the backing foam.\nSeveral attempts have been made to overcome this drawback. For example, special stabilizers for the PVC skin have been used. An intermediate layer between the skin and the backing foam to prevent the migration of plasticizers and additives (See EP No. 0,161,477) has also been employed.\nHowever, these measures which seriously restrict the choice of additives and/or lead to additional cost generating process steps do not result in a significant improvement.\nConsideration has been given to producing the solid skin and backing foam in a single step from a polyurethane system in the form of a so-called integral foam having a compact skin using the reaction injection molding process. In this case, however, reaction injection molding produces certain disadvantages, despite the saving of time. In such moldings, the skin and foam backing would have to be made of the same material. It would not be possible to make the skin, for example, from a light-stable, pigmented system and the foam backing from a simple, inexpensive system or to use systems having different mechanical properties for the skin and backing. In addition, it would be difficult to adjust the thickness of the skin as required within wide limits. The occasionally large and complicated undercuts of the moldings to be produced would also represent a certain problem in terms of mold design.\nAnother alternative for the production of a polyurethane skin is the in-mold coating process, in which the inside of the mold is first sprayed using a one-component or two-component spray or mixing head with enough material to form a layer a few tenths of a millimeter thick. However, considerable problems are created by the necessarily thin coating applied by a spray or mixing head. For example,the coating materials tend to run down vertical walls and mold undercuts are virtually impossible to achieve by this process.\nThe deep drawing of a polyurethane film is another method of forming a polyurethane skin. However, this process is subject to the same problems as the deep drawing of films of other materials. The deep-drawn skin has an irregular thickness (uneven stretching) and complicated mold geometries, such as undercuts are virtually impossible to achieve.\nA polyurethane skin could also be formed by the slush molding of a polyurethane thermoplast powder. This process is known in the case of PVC powders. However, mold temperatures above 180.degree. C. would be necessary and the energy costs would be very high."}
{"text": "1. Field of the Invention\nThe present invention is related to the field of telecommunications. More particularly, the present invention relates to removing slamming protection from a customer's account when the customer wants to change their preferred long distance provider.\nA portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent document or the patent disclosure, as it appears in the Patent and Trademark Office patent file or records, but otherwise reserves all copyright rights whatsoever.\n2. Acronyms\nAmeritech Service Order Negotiation (ASON)\nBilling Telephone Number (BTN)\nCarrier Identification Code (CIC)\nCarrier Request (CR)\nCustomer Account Record Exchange (CARE)\nElectronic 3-Way Call (E3WC)\nExtensible Markup Language (XML)\nFederal Communications Commission (FCC)\nFile Transport Protocol (FTP)\nGraphical User Interface (GUI)\nHypertext Markup Language (HTML)\nHypertext Transfer Protocol (HTTP)\nHypertext Transfer Protocol Secure (HTTPS)\nInteractive Voice Response (IVR)\nLocal Access and Transport Area (LATA)\nLocal Exchange Carrier (LEC)\nLocal Preferred Interexchange Carrier (LPIC)\nMultiple Virtual Storage (MVS)\nPreferred Interexchange Carrier (PIC)\nPublic Switched Telephone Network (PSTN)\nService Order Server Applications (SOSA)\nTelephone Number (TN)\nThird Party Verification (TPV)\nUniform Resource Locator (URL)\nWorking Telephone Number (WTN)\n3. Background and Material Information\nIn the past, if a potential customer consented to a change in their long distance carrier, then the long distance carrier would contact the customer's local telephone company to provision the change. Eventually, however, certain long distance companies began contacting local telephone companies and indicating that a potential customer had authorized a change in their long distance carrier, when in fact no such consent was granted by the customer. Ignorant of the falsehood, the customer's local telephone company would change the service as instructed by the carrier and the customer would receive service and incur charges from the newly provisioned long distance provider. This illegal practice eventually became known as slamming. In response to the advent of slamming, the Federal Communications Commission (FCC) adopted a requirement that the new service provider must first obtain third party verification (TPV) before a customer is switched to a new long distance carrier. The TPV is recorded and performed by an entity distinct from the long distance carrier and serves to verify that the customer consents to the change.\nThus, when a long distance carrier makes a sale to a customer, a TPV is performed and a recording is made of the telephone conversation. Then, the long distance provider submits a request to change the customer's long distance provider. If the customer has slamming protection on their account, then the local telephone company may not provision the change in long distance providers and notifies the long distance provider accordingly.\nAs a result, the long distance carrier will typically contact the customer and request that the customer contact their local telephone company and indicate that slamming protection be removed from their account. Alternatively, a three-way call may be attempted between the customer, the long distance provider, and the local telephone company. Often, a customer cannot be reached or does not have time to contact their local phone company. This is especially true when attempting to schedule a three-way call. In any event, the resulting frustration may result in lost sales for the carrier.\nTherefore, it would be advantageous to automate and streamline the process of removing slamming protection from a customer's account when the customer wants to change their preferred long distance provider."}
{"text": "Memory devices are typically provided as internal, semiconductor, integrated circuits in computers or other electronic devices. There are many different types of memory including volatile and non-volatile memory. Volatile memory can require power to maintain its data and includes random-access memory (RAM), dynamic random access memory (DRAM), and synchronous dynamic random access memory (SDRAM), among others. Non-volatile memory can provide persistent data by retaining stored data when not powered and can include NAND flash memory, NOR flash memory, read only memory (ROM), Electrically Erasable Programmable ROM (EEPROM), Erasable Programmable ROM (EPROM), and resistance variable memory such as phase change random access memory (PCRAM), resistive random access memory (RRAM), and magnetoresistive random access memory (MRAM), among others.\nMemory devices can be combined together to form a storage volume of a memory system such as a solid state drive (SSD). A solid state drive can include non-volatile memory (e.g., NAND flash memory and NOR flash memory), and/or can include volatile memory (e.g., DRAM and SRAM), among various other types of non-volatile and volatile memory.\nAn SSD can be used to replace hard disk drives as the main storage volume for a computer, as the solid state drive can have advantages over hard drives in terms of performance, size, weight, ruggedness, operating temperature range, and power consumption. For example, SSDs can have superior performance when compared to magnetic disk drives due to their lack of moving parts, which may avoid seek time, latency, and other electro-mechanical delays associated with magnetic disk drives.\nMemory is also utilized as volatile and non-volatile data storage for a wide range of electronic applications. Non-volatile memory may be used in, for example, personal computers, portable memory sticks, digital cameras, cellular telephones, portable music players such as MP3 players, movie players, and other electronic devices. Memory cells can be arranged into arrays, with the arrays being used in memory devices."}
{"text": "This invention relates to exercise equipment and, in particular, a treadmill having functionalities in addition to an adjustable running surface.\nWalking and running are two popular forms of exercise. They use substantially all muscles of the body and require very little mental concentration. Individuals who enjoy walking and running as forms of exercise also find that they are a therapeutic form of exercise as well. Therefore, treadmills have become popular because they may provide a programmable running surface with adjustable inclines and speeds in a very compact space indoors so that individuals may walk and/or run inside in a controlled environment during inclement weather. In addition, many runners prefer running on treadmills to running on pavement, which can be hard on their knees and other joints.\nHowever, as the world becomes more hectic, it has become increasingly difficult for individuals to maintain an exercise schedule, indoors or outdoors. Demands in the workplace include not only high stress, but also prolonged office hours. Office workers spend more time sitting at their desks, often with poor posture, and developing long-term back issues and undesired weight-gain, while, at the same time finding less and less time to exercise. With all the other demands of life such as family, friends, chores and hobbies, exercise often falls lower on the priority list of office workers. When this occurs, office workers may skip exercise as a part of their daily routine and, as a result, their health is compromised.\nU.S. Pat. No. 5,813,947 to Densmore, U.S. Pat. No. 5,984,839 to Corkum, U.S. Pat. No. 6,224,516 to Disch, and U.S. Pat. No. 6,283,896 to Grunfeld et al. are some other attempts that have been made to incorporate exercise equipment with computers, work stations, or other activity stations. However, with such distractions as computers, work stations or other activity stations, safety is compromised and a user of such inventions may be prone to falling down and sustaining injuries.\nTherefore, it would be desirable if the user could find a way to safely incorporate exercise into the workday, without sacrificing time that is already dedicated to the other demands of their lives."}
{"text": "There exists remote desktop software that allows a personal computer's desktop environment to be run on one system while being displayed on a separate client device. Often, this is done so that a human user may view and interact with the desktop of a remote computer. In this arrangement, the client device displays a remotely-presented, virtual, graphical window. It may access a single software application, such as a text editor, or computer game; or an entire computer “graphical desktop”, such as one sees when a common desktop computer has completely finished booting. Anything that the user accesses through the window displayed by the client device actually runs on the remote computer desktop or window manager. In contrast to both desktop applications as conventionally used and to conventional “Web applications”, in this “remote/virtual application window” (“remote window”) scenario the remote computer supplies almost all the “horsepower” for powering the software. The client uses rather generic logic to merely render the graphical user interface (along with any multi-media such as audio), take user input, take device input (such as geographical positioning coordinates, device orientation, or ambient lighting level), and relay all of it to the remote computer. The user may browse web content as usual, using hyperlinks, or any other standard navigational elements (generically termed “links”) inside these remote windows.\nThe same holds true for an automated user such as a search engine web crawler “bot” (also known as a “search engine web crawler”, or as a “web spider”, or as a “search bot”), in terms of browsing and interacting with this remote content. While rare, it is possible that a search bot would visit a computer desktop. More likely, however, software applications may be delivered in the same manner as a remote desktop, as described above, and those software applications may offer data, information, and content of interest to a search bot."}
{"text": "1. Field of the Invention\nThe present invention relates to rolls to be used in copying machines or printers, and particularly to conductive rolls to be provided around a photoconductor drum.\n2. Description of the Related Art\nIn a copying machine or a printer, an electrostatic latent image formed on a photoconductor drum is developed by toner (developer) fed from a developing roll. The toner image on the photoconductor drum is transferred to a transfer roll. The toner image is then transferred to a recording paper. Then, the recording paper is heated and pressed so that a picture image is formed on the recording paper. Such a roll as used in a copying machine or printer is generally has a structure in which an outer circumference of pipe-like core made of metal such as aluminum or iron is coated with rubber or resin material.\nWith reduction in hardness of such rolls, particularly, such conductive rolls including a charge roll, a developing roll, and a transfer roll which are used around a photoconductor drum, the roll is required to have high releasability to prevent the toner (developer) from adhering, for the purpose of ensuring nipping property between the roll and the photoconductor drum. Generally, each of the rolls such as conductive rolls adopts a structure to have a base material layer made of low hardness rubber material and provided on an outer circumference of a metal core, and a front surface layer made of resin material and provided on the base material layer.\nAs the rubber material constituting the base material layer of such a roll, low hardness rubber in which a softening agent such as oil is added to general-purpose rubber such as NBR or silicon rubber, or foamed rubber of polyurethane has been used. Especially silicon rubber has been used in view of stability in resistance value. As the resin material constituting the front surface layer, resin material have been used which is high in releasability, such as fluororesin, acrylate resin, urethane resin, or silicon resin.\nThe conductive rolls to be provided around a photoconductor drum are generally used with a difference in peripheral velocity between each conductive roller and the photoconductor drum due to the function of the conductive rollers. Especially, in the developing roller, taking the carrying quantity and developing property of toner into consideration, it is general to provide a peripheral velocity difference so that the peripheral velocity of the developing roll is about 1.1 times-1.5 times as large as the peripheral velocity of the photoconductor drum. Further, although it is general to particularly provide no peripheral velocity difference between the charge roll or the transfer roll and the photoconductor drum, there occurs actually a peripheral velocity difference in some degree due to variations in outer diameter of each roll per se.\nFurther, in the rolls other than the conductive rolls used around the photoconductor drum, for example, in the carrying rolls for carrying recording paper or in the fixing rolls, each roll is rotating in such a condition that a pair of rolls opposed to each other, that is, a roll in question and another roll which is a matter to be opposed to each other, are disposed to be in contact with each other. Accordingly, even if any peripheral velocity difference is not particularly provided between the rolls opposed to each other, it can be said that there exists, in fact, a certain amount of peripheral velocity difference between the thus opposed rolls.\nIf any peripheral velocity difference exists between a roll and a matter opposed to the roll, for example, in a case where a photoconductor drum is disposed as the matter opposed to the roll in question, it is apt to generate ripple-like deformation called wrinkles in the surface of the roll in question in the rotation driving. It is considered that such wrinkles are generated as follows. That is, as stated above, the roll is constituted by a base material layer of low hardness rubber material provided on the outer circumference of a metal core and a front surface layer of resin material provided on the base material layer. Accordingly, the soft base material layer is deformed by a peripheral velocity difference between the roll and a matter opposed to the roll in question, so that the deformation causes plastic deformation on the upper front surface layer.\nAs a method of preventing such ripple-like wrinkles from generating, it is considered that the front surface layer is made of a soft material, or the base material layer is made hard. However, when the front surface layer is made of a soft material, the lowering of the releasability is caused so that the function to be required originally cannot be satisfied. Further, the lowering of the releasability causes a disadvantage called filming due to adhesion of toner in the case where the roll in question is a developing roll. If the base material layer is made hard, on the contrary, it becomes difficult to ensure a sufficient nip width between the roll and an opposed matter such as a photoconductor drum. Accordingly, it is necessary to provide a countermeasure to heighten the accuracy of size of the roll in question per se so that a predetermined nip width can be obtained stably.\nFurther, in a case of a transfer roll which requires a certain nip width, generally, the roll in question could not help taking a soft roll structure at the sacrifice of the releasability. In addition to the problem of generation of ripple-like wrinkles, it is the existing condition that the mechanism has to be complex because a cleaning member is required to be provided to remove surplus toner adhering onto the roll surface.\nAn object of the present invention is to provide a roll which can secure necessary nip properties and releasability, and at the same time, the roll in question does not generate ripple-like wrinkles in the front surface layer even if there exists a peripheral velocity difference between the roll in question and any other opposed matter such as a photoconductor drum.\nIn order to achieve the object, according to the invention, there is provided a roll including: a base material layer made of low hardness rubber material and provided on an outer circumference of a metal core; and a front surface layer made of resin material and provided on the base material layer; wherein a stress relaxation layer is provided between the base material layer and the front surface layer, the stress relaxation layer being made of material which is at least higher in hardness than the base material layer and larger in stretch than the front surface layer.\nIn the roll of the invention, specifically, the base material layer is made of rubber material with low hardness not higher than 25 degrees in JIS-A hardness, the front surface layer is made of resin material not larger than 30% in stretch, and the stress relaxation layer is made of material having stress not lower than 5 MPa at 10% stretch.\nIn the roll of the invention, the stress relaxation layer is made of rubber material added with at least one kind of self-crosslinking resin such as blocked type self-crosslinking isocyanate resin, blocked type self-crosslinking epoxy resin, and blocked type self-crosslinking phenolic resin. The loading of the self-crosslinking resin to the rubber material is 20-150 parts by weight relative to 100 parts by weight of the rubber material.\nIn the roll according to the invention, the stress relaxation layer has a thickness in a range of from 5 to 50 xcexcm. In the roll according to the invention, preferably, the base material layer has a thickness not smaller than 3 mm, and the roll according to the invention is suitable as a conductive roll having roll hardness not higher than 60 degrees in Asker-C hardness at a load of 1 kg."}
{"text": "The present subject matter relates generally to systems and methods for defrosting eutectic plates of a refrigeration system.\nRefrigeration plates for maintaining a low temperature in a storage area are well known. For example, trucks with refrigerated storage areas are commonly used for transporting frozen or perishable foods. In a typical storage area, a refrigeration unit comprises a plurality of metal plates. Each metal plate has a cavity filled with an amount of freezable fluid or eutectic solution and a conduit immersed in the eutectic solution. The conduit winds back and forth through the interior of the plate, with the two ends of the conduit communicating with the outside of the plate. An exemplary refrigeration plate is described in U.S. Pat. No. 3,756,037, the disclosure of which is hereby incorporated herein by reference.\nThe ends of the conduit are connected in line with a refrigeration circuit including, among other things, a condenser, a compressor, and a volatile refrigerant that flows through the circuit. An exemplary refrigeration circuit is described in U.S. Pat. No. 4,043,144, the disclosure of which is hereby incorporated herein by reference. A “freezing cycle” is implemented to freeze the eutectic solution during a period of non-use of the refrigerated storage area. Once the eutectic solution has frozen and the plates have reached a sufficiently low temperature, the freezing cycle is terminated and the storage area may be returned to use during an operating period. During the operating period, air is blown across the plates to cool the storage area while the truck is used to deliver the cargo to its destination.\nOver the course of the operating period, frost may accumulate on the surface of the plates due to moisture in the air. Since frost acts as a thermal insulator, the presence of frost on the plates and between adjacent plates reduces their cooling efficiency, so the plates must be periodically defrosted.\nThere are several known methods for defrosting the plates. For example, U.S. Pat. No. 3,727,422 describes a method of applying water to the surfaces of the plates to remove frost. In another system, which is described in U.S. Pat. No. 4,043,144, a defrost coil is positioned against the outer surfaces of the plates and hot gas is run therethrough to defrost the plates. U.S. Pat. No. 6,595,019 describes yet another system, which employs vertical scrapers that are moved across the outer surfaces of the plates to remove frost. All of these patents are hereby incorporated herein by reference.\nAll of these methods have certain disadvantages. For example, using water to defrost the plates requires an operator to connect a water hose to a water dispersal or spray system to begin the defrosting process by running relatively warm water over the frost. Then the water falling off of the plates has to be collected and disposed of. And finally the operator has to disconnect the hose when the defrosting process is completed.\nIn view of the foregoing, there is a need for a defrost system and method which are not only effective, but which also eliminate the need for human intervention for defrosting control."}
{"text": "Reactive Ion Etch (RIE) is used to form shapes on work pieces such as semiconductor wafers. In a typical RIE process, radio frequency (RF) or microwave power is used to excite a gas to form a plasma. The plasma facilitates etching the shapes into the work piece. A typical use of the RIE process is in etching very fine holes, typically minimum dimension sized holes, through a layer, e.g., an insulating layer, to an underlying layer, commonly referred to as \"opening contacts\" or \"opening vias\" through the layer. RIE is also used to open minimum width trenches through one or more work piece layers. Typically, these minimum dimension shapes are deeper than at least one of their surface dimensions and, therefore, have a high aspect ratio.\nFIG. 1A is a schematic diagram of a conventional diode RIE system 100, such as an MXP chamber of the AME 5X00 series, manufactured by Applied Materials Corp. The RIE system includes an etching chamber 2. A window 4 is formed in one portion of etching chamber 2 for viewing the RIE process from outside etching chamber 2. A slot 6 is formed in another portion of etching chamber 2 through which a work piece such as a semiconductor wafer may be inserted for etching. A pump 8 is connected to etching chamber 2 for creating a vacuum in etching chamber 2. Magnetic coils 10a, 10b, 12a, and 12b are disposed about the perimeter of etching chamber 2. Each magnetic coil is offset by about 90 degrees from the magnetic coils on either side. For example, magnetic coil 10a is offset by about 90 degrees from magnetic coils 12a and 12b, respectively. Magnetic coil 10a faces magnetic coil 10b across etching chamber 2 and, similarly, magnetic coil 12a faces magnetic coil 12b across etching chamber 2.\nMagnetic coils 10a, 10b, 12a, and 12b are typically driven by an AC power source (not shown) providing sinusoidal functions. In particular, magnetic coils 10a and 10b are driven by a first sinusoidal function provided by the AC power source, and magnetic coils 12a and 12b are driven by a second sinusoidal function offset by approximately 90 degrees with respect to the first sinusoidal function. As such, at a first phase corresponding to a maximum of the first sinusoidal function, a magnetic field passes between magnetic coils 10a and 10b across etching chamber 2. Similarly, at a second phase corresponding to a maximum of the second sinusoidal function, a magnetic field passes between magnetic coils 12a and 12b. In this way, magnetic coils 10a, 10b, 12a, and 12b produce a magnetic field which rotates over time with respect to etching chamber 2. The magnetic field produced by magnetic coils 10a, 10b, 12a, and 12b typically rotates at a rate of about 0.2 Hertz (Hz).\nFIG. 1B illustrates a cross-sectional view of the conventional diode RIE system 100 of FIG. 1A. Gas is injected into etching chamber 2 through shower head 24 and is pumped out of etching chamber 2 by pump 8. A semiconductor wafer 14 has been inserted through slot 6 (shown in FIG. 1A), dropped onto a work piece chuck 16, and clamped to a work piece chuck electrode 18 formed as a portion of work piece chuck 16. Work piece chuck electrode 18 is driven by a power supply 20. Power supply 20 creates a radio frequency (RF) field through work piece chuck electrode 18 which, in turn, produces a plasma 22. Plasma 22 may diffuse slightly along the sides of work piece chuck 18. A rotating electromagnetic field 26 is produced by magnetic coils 10a, 10b, 12a, and 12b, as described with reference to FIG. 1A. Specifically, FIG. 1B illustrates electromagnetic field 26 at a time corresponding to the first phase described above, such that rotating electromagnetic field 26 passes between magnetic coils 10a and 10b across etching chamber 2.\nBecause electron plasma 22 is generated by AC power source 20 through work piece chuck electrode 18, electron plasma 22 is capacitively driven. As such, the plasma efficiency of conventional diode RIE system 100 decreases as the density of plasma 22 increases. Therefore, for the conventional diode RIE system to operate efficiently, the density of plasma 22 is low. Such an RIE system 100 often results, however, in less-than-desirable etch rates, plasma non-uniformities, and ionic charging.\nEffects from ionic charging, known to those skilled in the art as \"aspect ratio charging effects,\" typically occur at the bottom of high aspect ratio structures, especially if some part of the structure is an insulator. These aspect ratio charging effects include oxide damage, device damage, threshold voltage shifts, polysilicon notching, and reduction of ion current at the bottom of trenches and at the bottom of vias that impede or even stop etching (known as RIE lag).\nAs line widths have decreased and aspect ratios of etch structures have increased, aspect ratio charging effects have become more problematic. These charging effects are due, at least in part, to the fact that, in the conventional RIE system, electrons diffuse out of the plasma due to high thermal energy and densities while positive ions are extracted. Thus, positive ions are accelerated to various locations on the structure and electrons spread and are collected near mask surfaces until a potential develops which equalizes these two particle fluxes. There remains a need, therefore, to eliminate aspect ratio charging effects."}
{"text": "The present invention relates to a bottom structure of a paper box for storage of liquid, and more particularly, relates to a paper box for storage of liquids such as milk, juice, sake, soy sauce and the like.\nRecently, various paper boxes made of thermoplastic synthetic resin laminate plate-like paper have been used instead of conventional metal cans or glass bottles as containers for storing liquid such as milk, juice or sake. The reason for this is that the paper boxes have the following advantages;\n(a) Since the paper boxes have generally a regular parallelepiped shape, any dead space can be removed and they may be laid one on another in carriage or storage; PA1 (b) The paper boxes are light in weight and thin in thickness in comparison with glass bottles; PA1 (c) The value of goods may be enhanced by applying decorative printing on their outside surface; PA1 (d) The paper boxes may readily be collapsed flat after consuming contents therein; PA1 (e) The paper boxes may be disposed of by burning after use; and PA1 (f) The content contained therein may be protected securely.\nConventionally, as such a paper box, there has been a container developed, manufactured for sale by Ex-Cello-O Corporation (Japanese Utility Model Publication No. 4661/71) for milk containers. As shown in FIG. 1, its bottom structure has a plate-like paper structure as indicated by folding lines, i.e., dot-and-dash lines and folding-back lines, i.e., by two-dot-one-dash lines. After a connecting portion G formed along one edge of the plate-like paper is connected with corresponding portion opposite thereto to form a rectangular tube as shown in FIG. 2A, the paper is folded along the folding lines and the folding-back lines as shown in FIGS. 2B to 2D to thereby form the bottom structure. In the container disclosed in the above-mentioned Japanese Utility Model Publication No. 4661/71, one bottom plate A among the two bottom plates A and B forming outermost bottom surface is inserted between the confronting bottom plate B and the folded inside portions as shown in FIGS. 2A to 2D. Since the edges of the folded portions C and D, E and F are intimately contacted with each other by a precise arrangement for folding, the container has the above-described advantages that the liquid contained therein will not leak, however there has been a fear that, when the end of the bottom plate is inserted into an extremely narrow gap between bottom plate B and folded portions D, E as shown in FIG. 2C, the end of the bottom plate A may be folded back midway in the process.\nIn order to overcome these defects, a paper box having the end portion of the bottom plate A removed as shown in FIG. 3 has been manufactured and available in the market. However, the bottom structure of the paper box is made up of a half-size bottom plate which is overlapped in part only by the other half-size bottom plate when folded, and, in association with the folding-back portions C, D, E and F, each edge thereof does not always abut precisely with the others unlike the preferred paper box illustrated by FIG. 4A, but poor overlapping is likely to be caused as in the case of sub-standard boxes illustrated by FIG. 4C in which case liquid contained therein is threatened to leak through gaps therebetween and the boxes of such structure have eventually led to a deteriorated quality of box."}
{"text": "Immediate processing of data is very important when a device within a computing system runs a real time application. Such real time applications may involve the processing of video or audio. The device may require adequate bandwidth from one or more sources in order to effectively display images or reproduce audio. The device may require access to data stored in a memory of the computing system.\nSuch computing systems may utilize a memory controller to access data stored in memory. The memory controller may employ a device called a memory arbiter that handles requests from one or more devices requiring access to data stored in memory. These devices may communicate to the memory using one or more data buses resident in the computing system.\nOften, a device running a real time application may not be able to receive data from main memory at a required rate. One or more other devices performing non-real time applications may consume bandwidth required by the device running the real time application. For example, a digital camcorder may require transmission of data from main memory at a sufficient rate in order to properly display video on a monitor. When other devices compete for memory resources, access to memory may not be provided to the most critical applications such as those which involve real time processing of data.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "The thin film solar cell which used CuInSe2 (CIS based thin film) which is a semiconductor thin film of chalcopyrite structure composed of a group Ib element, a group IIIb element, and a group VIb element, or Cu(In, Ga)Se2 (CIGS based thin film) which dissolved Ga to the CuInSe2 for an optical absorption layer shows high energy conversion efficiency, and has advantage that there is little degradation of the efficiency by light irradiation etc.\nHowever, film formation by 550 degrees C. from a viewpoint of degradation of film quality and increase of leakage current is general, in formation of the CIS based thin film which is the semiconductor thin film of chalcopyrite structure, or the CIGS based thin film which dissolves Ga to CIS based thin film. When it forms at low temperature rather than 550 degrees C., it has been considered that particle diameter is small composed and dark current characteristics degrade, conventionally. In addition, the heat-resistant limitation of an integrated circuit is about 400 degrees C.\nIt is already disclosed about a photoelectric conversion device which used a compound semiconductor thin film of chalcopyrite structure to reduce dark current substantially, and a fabrication method for the such photoelectric conversion device (for example, refer to Patent Literature 1 and Patent Literature 2).\nAlso, it is already disclosed about a formation method by a selenium process of a high quality CIGS based thin film (for example, refer to Patent Literature 3 and Patent Literature 4).\nOn the other hand, it is already disclosed about: a solid state imaging element characterized by forming a switching device by a thin film transistor on a substrate and composing to laminate a sensor area by an amorphous semiconductor layer via the picture element electrode connected to the above-mentioned switching device; or a solid state imaging element composed to form by an insulating substrate as the above-mentioned substrate (for example, refer to Patent Literature 5).\nIn the solid state imaging element disclosed in Patent Literature 4, since the amorphous semiconductor layer is applied into the photosensing area, a photoelectric conversion wavelength is mainly a visible light wavelength region.\nIn such a conventional solid state imaging element, since a low electric field is applied to a photoelectric conversion film and an electric charge is detected, the photoelectric conversion film itself has no multiplication function.\nCitation List\nPatent Literature 1: Japanese Patent Application Laying-Open Publication No. 2007-123720\nPatent Literature 2: Japanese Patent Application Laying-Open Publication No. 2007-123721\nPatent Literature 3: U.S. Pat. No. 5,436,204\nPatent Literature 4: U.S. Pat. No. 5,441,897\nPatent Literature 5: Japanese Patent Application Laying-Open Publication No. 2001-144279"}
{"text": "The polyphenylene ether resins are a family of engineering thermoplastics that are well known to the polymer art. These polymers may be made by a variety of catalytic and non-catalytic processes from the corresponding phenols or reactive derivatives thereof. By way of illustration, certain of the polyphenylene ethers are disclosed in Hay, U.S. Pat. Nos. 3,306,874 and 3,306,875, and in Stamatoff, U.S. Pat. Nos. 3,257,357 and 3,257,358. In the Hay patents, the polyphenylene ethers are prepared by an oxidative coupling reaction comprising passing an oxygen-containing gas through a reaction solution of a phenol and a metal-amine complex catalyst. Other disclosures relating to processes for preparing polyphenylene ether resins, including graft copolymers of polyphenylene ethers with styrene type compounds, are found in Fox, U.S. Pat. No. 3,356,761; Sumitomo, U.K. Pat. No. 1,291,609; Bussink et al. U.S. Pat. No. 3,337,499; Blanchard et al. U.S. Pat. No. 3,219,626; Laakso et al. U.S. Pat. No. 3,342,892; Borman, U.S. Pat. No. 3,344,166; Hori et al. U.S. Pat. No. 3,384,619; Faurote et al. U.S. Pat. No. 3,440,217; and disclosures relating to metal based catalysts which do not include amines, are known from patents such as Wieden et al., U.S. Pat. No. 3,442,885 (copper-amidines); Nakashio et al., U.S. Pat. No. 3,573,257 (metal-alcoholate or -phenolate); Kobayashi et al. U.S. Pat. No. 3,455,880 (cobalt chelates); and the like. In the Stamatoff patents, the polyphenylene ethers are produced by reacting the corresponding phenolate ion with an initiator, such as peroxy acid salt, an acid peroxide, a hypohalite, and the like, in the presence of a complexing agent. Disclosures relating to non-catalytic processes, such as oxidation with lead dioxide, silver oxide, etc. are described in Price et al., U.S. Pat. No. 3,382,212. Cizek, U.S. Pat. No. 3,383,435 discloses polyphenylene ether-styrene resin compositions. All of the above-mentioned disclosures are incorporated by reference.\nThe term \"polystyrene resin\" includes polymers and copolymers of styrene, alpha methyl styrene, chlorostyrene, and the like.\nThe term \"EPDM\" includes rubbery interpolymers of a mixture of mono-olefins and a polyene. Preferred types are those rubbery interpolymers of ethylene, an alpha-olefin, and a polyene. Rubbery interpolymers of ethylene, propylene, and a polyene are especially preferred.\nThe term \"EPDM-silicone rubber\" includes graft copolymers of EPDM with silicone rubber. An example of an EPDM-silicone rubber employed in the present invention is that commercially manufactured and sold by Shinetsu Chemical Industry, under the name SEP-172U.\nIn the prior art, rubber-modified styrene resins have been admixed with polyphenylene ether resins to form compositions that have modified properties. The Cizek patent, U.S. Pat. No. 3,383,435, discloses rubber-modified styrene resin-polyphenylene ether resin compositions wherein the rubber component is of the unsaturated type such as polymers and copolymers of butadiene. The physical properties of these compositions are such that it appears that many of the properties of the styrene resins have been upgraded, while the moldability of the polyphenylene ethers are improved.\nNakashio et al. U.S. Pat. No. 3,658,945 discloses that from 0.5 to 15% by weight of an EPDM-modified styrene resin may be used to upgrade the impact strength of polyethylene ether resins. In Cooper et al., U.S. Pat. No. 3,943,191 it is disclosed that when the highly unsaturated rubber used in compositions of the type disclosed by Cizek, is replaced with EPDM rubber that has a low degree of residual unsaturation, the thermal oxidative stability and color stability are improved. The EPDM rubber in the Cooper et al. compositions is comprised substantially of particles in the range of 3-8 microns.\nThe impact strength of the Cooper et al. compositions is superior to that of a polypropylene ether resin alone or that of similar compositions comprised of unmodified polystyrene; however, the impact strength of the Cooper et al. compositions is inferior to that of similar compositions comprised of polystyrene modified with polybutadiene rubber, such as a composition known as FG-834 available from Foster-Grant Co. As is disclosed in U.S. Pat. No. 3,981,841, the impact strength of the Cooper et al. compositions can be improved by incorporating therein impact modifiers such as an emulsion-grafted EPDM polystyrene copolymer. U.S. Pat. No. 4,152,316 incorporated herein by reference, discloses that a composition of a polyphenylene ether resin and an alkenyl aromatic resin modified with an EPDM rubber comprised of particles having a median diameter less than about two microns, preferably about 0.5 to 1.5 microns, is a very useful thermoplastic molding material having good thermal oxidative stability and good room temperature impact strength but inferior low temperature impact strength.\nIn U.S. Pat. No. 4,102,850 it is disclosed that the addition of small amounts of mineral oil to the polymerizing mixture of styrene and EPDM rubber produces EPDM-modified polystyrene which yields blends with polyphenylene oxide having significantly better low-temperature impact strength than blends made from EPDM-polystyrene made without the mineral oil.\nIn U.S. Pat. No. 3,737,479 it is disclosed that the addition of a silicone such as a polyorganosiloxane, which is fluid, to polyphenylene oxide or to polyphenylene oxide-polystyrene blends improves Gardner impact strength but not Izod impact strength (see Col. 1, lines 54-57 of U.S. Pat. No. 3,737,479).\nIt has now been found when small amounts of a graft copolymer of EPDM with silicone rubber is blended with a polyphenylene oxide resin, and optionally with a polystyrene resin, blends can be obtained which are substantially improved in ductility and Izod impact strength both at room temperature and at low temperatures."}
{"text": "1. Field of the Invention\nThe present invention relates to a fuel cell system.\nPriority is claimed on Japanese Patent Application No. 2005-177592, filed Jun. 17, 2005, the content of which is incorporated herein by reference.\n2. Description of the Related Art\nIn a known fuel cell system having a plurality of sensors for detecting operation states of a fuel cell, when a specific sensor (e.g., an air flow rate sensor, a hydrogen inlet pressure sensor, or a cooling water inlet pressure sensor) has an abnormal state, a substitutional value with respect to a signal value output from the sensor in the abnormal state is computed based on a signal value output from another sensor (e.g., an air temperature sensor, a hydrogen outlet pressure sensor, or a cooling water outlet pressure sensor) in a normal state, and the operation of the fuel cell is continued in accordance with the computed substitutional value (see, for example, Japanese Unexamined Patent Application, First Publication No. 2004-179127).\nIn the fuel cell system relating to the above conventional technique, the fuel cell is operated in accordance with signal values output from the plurality of sensors, and only a substitutional value for the signal value output from a sensor in the abnormal state is computed based on a signal value output from another sensor in the normal state. Therefore, if the sensor necessary for computing the substitutional value enters an abnormal state, it is difficult to continue the operation of the fuel cell. For example, when the fuel cell system is provided for supplying electric power to a driving source of a vehicle, it may be difficult to run the vehicle.\nIn addition, if redundant or excessive sensors are provided for computing a substitutional value for a target sensor, the structure of the fuel cell system is complicated, thereby increasing the cost necessary for implementing the structure."}
{"text": "This invention relates in general to the inspection of surfaces to detect anomalies, and in particular, to an improved system that illuminates the surface inspected at the plurality of spots simultaneously for anomaly detection.\nConventional optical inspection methods employing scanning techniques typically causes a single spot on the surface inspected to be illuminated where the spot is scanned over the entire surface for anomaly detection. For improved signal-to-noise ratio caused by background scattering, the size of the illuminated spot has been continually reduced. This means that the amount of time required for the spot to scan over the entire surface is increased which is undesirable.\nOne solution to the above dilemma is proposed in U.S. Pat. No. 6,208,411, which is incorporated herein by reference in its entirety. This patent proposes a massively parallel inspection and imaging system which illuminates the surface at a plurality of spots where scattered light from the spots are imaged onto corresponding detectors in a detector array.\nWhile the system in U.S. Pat. No. 6,208,411 provides a major enhancement in the total inspection throughput, it may be further improved for enhanced performance in certain applications. It is, therefore, desirable to provide an improved multi-spot inspection and imaging system with enhanced characteristics."}
{"text": "With the development of smart electronic devices, there has been an explosive growth in various applications.\nCurrently, most applications have a function of automatically pushing messages. For a user, some of the pushed messages are useful, while some message pushed from some applications are simply interferences. In order to satisfy the user's requirements, the user needs to configure the applications individually to enable the message pushing function of some applications and disable the message pushing function of the other applications. As the number of the applications increases, such configuration process becomes too troublesome.\nIn addition, most electronic devices have many applications installed thereon. When a user is using an electronic device, an application may synchronize data via a network in background even if it has been deactivated. For example, application clients such as weather or news in a mobile phone may send or receive relevant data anytime for synchronization, such that a large amount of data traffic will be consumed without the user's awareness. In this case, the electronic device may be demanding on network traffic. In order to prevent an application from consuming traffic in background, the user typically needs to disable the synchronization setting for the application manually. Alternatively, the user may manage applications by installing application management software (e.g., safeguard) to configure an application to be disconnected from the network based on a timer or prompt a traffic limit. However, these configuration or management schemes require troublesome operations."}
{"text": "This invention relates to a position detecting device employed in an automatic positioning system.\nFIG. 1 shows a position detecting system using limit switches which is the most typical of the conventional position detecting systems. For simplification in description, the drive system is not shown in FIG. 1.\nAs shown in FIG. 1, a moving table 30 has a protrusion 30a which is used to detect the positions of limit switches 31a, 31b and 31c. As the moving table 30 is moved, the protrusion 30a operates the actuators 310 of the limit switches 31a, 31b and 31c, so that switch signals q, r and s are applied to a controller 32.\nThe position detecting limit switches 31c, 31b and 31a are set at distances L.sub.1, L.sub.2 and L.sub.3 from the original point O, respectively.\nThe conventional device thus constructed operates as follows. When the moving table 30 is moved in the direction of the arrow to the limit switch 31c at the distance L.sub.1 from the original point O, the actuator 310 of the limit switch 31c is operated, and the switch signal s is applied to the controller 32. Similarly when the table 30 reaches the limit switch 31b at the distance L.sub.2, the limit switch 31b is operated and the switch signal q is applied to the controller 32. When the table 30 reaches the limit switch 31a at the distance L.sub.3, the limit switch 31a is operated and the switch signal r is applied to the controller 32. In response to the switch signals q, r and s, the controller 32 operates to control the driving, stopping or moving speed of the moving table 30.\nEven when the controller 32 provides an instruction signal to stop the table 30, the table 30 will overrun a relatively long distance before it is actually stopped, because of the moving table's inertia and the delayed braking operation of the brake member. Accordingly, in this system, the limit switches 31a, 31b and 31c are located somewhat before the desired stopping points.\nThe overrun distance L can be approximately represented by the following expression: ##EQU1## where M is the mass of the moving part of the system, K is the rigidity of the mass in the feed direction, and V is the feed speed. Thus, it is difficult to eliminate the overrun distance L. Therefore, in the conventional system, in order to minimize the overrun distance it is necessary not only to set the limit switches 31a, 31b and 31c before the stop points but also to decrease the feed speed of the moving table 30 at the detection points.\nHowever, the overrun distance varies with the speed of the moving table and it is therefore difficult to determine the proper lead time. In addition, when the table 30 is moved in the opposite direction (or returned), the positions of the limit switches 31a, 31b and 31c will now be somewhat behind the desired stopping points and must be adjusted again.\nPosition detecting sensors employed in the conventional positioning system can be classified into mechanical types such as a jib, a cam, a link or a gear device; electrical types such as a microswitch or a contactless switch; pneumatic types such as a pneumatic microswitch; or optical types such as a photoelectric switch, a photoelectric linear scale or an optical wave interferometer.\nIn the case of the mechanical type position detecting sensor, the drive sources (such as a hydraulic drive source and/or an electric motor) are, in general, stopped each time a position is detected. Therefore, the mechanical type position detecting sensor is not suitable for automatic positioning. Even if it were designed so as to be suitable for automatic positioning, it would need considerably intricate mechanical elements. In addition, because of the inertial forces and drive forces of the mechanical elements, other mechanical elements may be deformed or worn. Therefore, with the mechanical type sensor, it is difficult to improve the positioning accuracy, and it is considerably difficult to change the design or to move the detection points as discussed above in the case of reverse direction movement.\nA microswitch or a contactless switch has been extensively employed as the electrical type position detecting sensor. The electrical type position detecting sensor is usually a simple on-off type digital detector. With such a position detecting sensor, however, a detection accuracy higher than that defined by the switch configuration cannot be obtained; that is, the detection accuracy is limited by the number and location of switches and cannot be improved. Furthermore, the number of sensors which can be installed over the range of movement of the moving table is limited, and it takes a great deal of time and labor to change the design and to adjust the detection points as discussed above.\nIn the case of the pneumatic type position detecting sensor, once a positioning point is determined it is considerably difficult to change the positioning point. Also, since an external force is repeatedly imparted to the position detecting element, it is rather difficult to maintain the accuracy unchanged.\nIn general, the optical type position detecting sensor is low in accuracy. Some high accuracy optical type position detecting sensors are available, but they are considerably expensive. Since the detection accuracy of the sensor is greatly dependent on the selected light source, the adjustment takes time and labor, and it is also difficult to change the positioning point.\nThus, the following characteristics are common in conventional devices:\n(1) It takes a great deal of time and labor to set the detection positions;\n(2) The number of detection positions is limited by the physical dimensions of the position detecting sensor used, and cannot be increased. The detection accuracy is also limited; and\n(3) After the detection positions are set, it takes a lot of time and labor to change the detection positions.\nA conventional positioning device comprises: a position detecting sensor; and a control unit for driving and controlling the device in response to the output signal of the sensor. The position detecting sensor suffers from the various drawbacks as described above."}
{"text": "The present invention generally relates to an apparatus and a method for filling holes in a substrate and more particularly, relates to an apparatus and a method for filling holes in electronic substrates that have high aspect ratios of at least 5:1 in a process for forming vias and interconnects.\nIn the electronics packaging industry, there is often a need to fill holes in various substrates for forming vias and interconnects. The hole filling process is affected by a number of processing and material parameters, for instance, the diameter of the via hole, the depth of the hole and the type of the substrate material, etc.. Typically, via holes are through holes and are used to electrically connect a top surface and a bottom surface of a substrate, the filler material utilized in forming the vias must be electrically conductive. A number of techniques have been used to perform the via hole filling process. Depending on the type of the substrate, the techniques may include electroplating, electroless plating, solder paste screen printing and conductive paste screen printing.\nA fundamental processing problem arises when the depth-to-width ratio, or the aspect ratio grows. The problem gets more complex when the hole diameter is made aggressively small, i.e., 125 xcexcm or less. For instance, for holes having depth-to-width aspect ratios of 5:1 or larger, it is no longer possible to use a conventional technique such as paste screening to fill the holes. Even the more advanced plating techniques cannot be used when diameters further decrease and aspect ratios further increase.\nAn attempt to fill via holes that have 17:1 aspect ratio with conductive epoxies by a screen printing method was found ineffective, even when the screening process is conducted from both the top and the bottom side of the substrate. One other attempt to fill a glass substrate that has via holes of 17:1 aspect ratio was carried out by utilizing a conductive material that has extremely low viscosity, i.e., an eutectic solder that has a viscosity of only 2 centipoise in its molten state, approaching that of water. Even at such low viscosity, known techniques for filling the high aspect ratio via holes were found ineffective. For instance, the conventional techniques tried include an injection molded solder technique with only pressure utilized. The injection molded solder technique was not able to fill deep vias holes that have very small diameters, i.e., smaller than 25 xcexcm.\nA second technique of vacuum injection molding was also tried. A graphic illustration of the vacuum injection molding process is shown in FIGS. 1Axcx9c1D. The vacuum injection molding process was disclosed in a co-pending application that was assigned to the common assignee of the present invention under Ser. No. 08/518,874 and now U.S. Pat. No. 6,231,333 B1. The vacuum injection molding method utilizes a pressure differential formed between either ambient and vacuum or positive pressure and vacuum. The pressure and the vacuum are both supplied on the same surface of a substrate 10, as shown in FIG. 1A. The process is carried out by utilizing a shallow vacuum link 12 that allows a continual evacuation of air from via holes 20 that have a large aspect ratio such as 5:1. The vacuum link 12 must be sufficiently shallow such that the surface tension of molten solder prevents cross-leaking during the operation. Such a shallow link 12 effectively choke a significant part of the full pressure differential and thus producing only partial filling of via holes 20 that have high aspect ratios.\nAs shown in FIG. 1A, an injection head 14 which includes spaced apart vacuum slot 16 and injection slot 22 is positioned on top of a mold plate, or substrate 10 in fluid communication with the mold cavities, or via holes 20 contained therein. Relative axial sliding is effected between the injection head 14 and the mold plate 10 for sequentially evacuating gas from the mold cavities 20 using a continuous vacuum and injecting into the evacuated mold cavities 20 a liquid fed from a continuous source (not shown). The sliding of the injection head 14 over the mold plate 10, as shown in FIGS. 1B and 1C, automatically provides self valving for sequentially evacuating and filling the mold cavities 20 from the same side, i.e., the top side 18 of the mold plate 10. In a preferred embodiment, the vacuum slot 16 and the injection slot 22 are linked together at the mold plate 10 so that surface tension of the liquid restrains flow of the liquid (not shown) from the injection slot 22 to the vacuum slot 16 while allowing gas flow thereinbetween for effecting vacuum in the mold cavities 20.\nAs shown in FIG. 1A, at the start of the process, i.e., before scanning begins, the mold cavities 20 are empty. As the scanning process begins, as shown in FIG. 1B, some cavities 24 are evacuated. As the scanning of the injection head 14 continues, some cavities 26 are filled with the liquid. In the final step of the process, as shown in FIG. 1D, the vacuum supply to the vacuum slot 16 is turned off while the scanning of the injection head 14 is completed, i.e., all the cavities 20, 24 and 26 are filled with liquid.\nThe vacuum injection molding method illustrated in FIGS. 1Axcx9c1D applies a vacuum and a liquid injection on the same surface of the substrate. To enable the operation, a shallow vacuum link must be used to allow the continual evacuation of air from mold cavities that have a larger aspect ratio such as 5:1. The drawback of the process is that the vacuum link must be sufficiently shallow such that the surface tension of molten solder prevents cross-leaking during the operation. The shallow link chokes off a significant part of the full pressure differential and thus only produces partial filling of the mold cavities, or the via holes that have high aspect ratios.\nIt is therefore an object of the present invention to provide an apparatus for filling a liquid in high aspect ratio holes that does not have the drawbacks or shortcomings of the conventional apparatus.\nIt is another object of the present invention to provide an apparatus for filling a liquid in high aspect ratio holes in a substrate that applies a vacuum and a liquid injection on the opposite sides of a substrate.\nIt is a further object of the present invention to provide an apparatus for filling a liquid in high aspect ratio holes in a substrate by utilizing a filler plate and a vacuum plate sandwiching a substrate thereinbetween such that vacuum and liquid injection are applied on the opposite sides of the substrate.\nIt is another further object of the present invention to provide an apparatus for filling a molten solder in high aspect ratio via holes in an electronic substrate which is effective in filling holes that have aspect ratios of at least 5:1.\nIt is still another object of the present invention to provide an apparatus for filling a molten solder in high aspect ratio via holes in an electronic substrate wherein a filler plate and a vacuum plate are connected together and are further equipped with a veneer adjustment means such that a relative position of the two plates can be adjusted corresponding to the viscosity of the liquid injected.\nIt is yet another object of the present invention to provide an apparatus for filling a liquid in high aspect ratio holes in an electronic substrate by pulling vacuum and injecting molten solder simultaneously from and into the same holes when the liquid injected has high viscosity.\nIt is still another further object of the present invention to provide a method for filling high aspect ratio via holes in an electronic substrate by pulling a vacuum from via holes on a first side of the substrate and injecting liquid on the opposite side of the substrate into the via holes.\nIt is yet another further object of the present invention to provide a method for forming interconnects in a substrate for a display panel such that a voltage potential can be placed on pixel display elements formed on the display panel.\nIn accordance with the present invention, an apparatus and a method for filling high aspect ratio via holes in electronic substrates are provided.\nIn a preferred embodiment, an apparatus for filling a liquid in high aspect ratio holes in a substrate is provided which includes a filler plate adapted for receiving an injection head therein, the filler plate has a bottom surface adapted for receiving an injection slot provided on the injection head, a vacuum plate adapted for receiving a vacuum slot in a top surface, the vacuum slot is in fluid communication with a vacuum source, and a connection means for connecting the filler plate and the vacuum plate together in a face-to-face, spaced-apart relationship forming a gap thereinbetween with the bottom surface of the filler plate positioned parallel and opposite to the top surface of the vacuum plate, the connection means further includes an adjustment means for adjusting a relative axial position of the filler plate to the vacuum plate, the gap formed between the filler plate and the vacuum plate is sufficiently large to slidingly engaging a substrate therein and to form abutting contacts between the substrate and the bottom surface of the filler plate and between the substrate and the top surface of the vacuum plate so that holes in the substrate are first evacuated by the vacuum slot and then filled with a liquid by the injection slot in the filler plate.\nIn the apparatus for filling a liquid in high aspect ratio holes in a substrate, the adjustment means may include a veneer adjustment for fixing a relative axial position of the filler plate to the vacuum plate. The adjustment means may further include a veneer adjustment for transforming a circumferential displacement into a linear axial displacement. The filler plate may be positioned by the adjustment means to lag behind the vacuum plate such that the injection slot lags behind the vacuum slot. The filler plate may be positioned by the adjustment means at the same axial position of the vacuum plate such that the injection slot and the vacuum slot are in fluid communication with a hole in the substrate simultaneously.\nIn the apparatus for filling a liquid in high aspect ratio holes in a substrate, the injection slot may have a width sufficiently large to cover substantially all the high aspect ratio holes in the substrate. The injection head may be in fluid communication with a liquid reservoir for feeding a liquid through the injection slot. The injection slot may be adapted for feeding a high electrical conductivity liquid solder therethrough. The connection means may further include an adjustment means of a threaded veneer mechanically attached to the filler plate and at least one engagement pin connecting the threaded veneer to the vacuum plate. The high aspect ratio holes in the substrate are selected from the group consisting of via holes, deep through holes and deep trenches.\nIn the apparatus for filling a liquid in high aspect ratio holes in an electronic substrate, the high aspect ratio holes may have an aspect ratio of at least 5:1, or an aspect ratio of at least 20:1. The high aspect ratio holes may have a diameter not smaller than 5 xcexcm. The electronic substrate may be adapted for making sliding movement between the filler plate and the vacuum plate at a speed of between about 25 mm/min and about 250 mm/min. The apparatus may further include a drive means slidingly engaging the substrate between the filler plate and the vacuum plate.\nThe present invention is further directed to a method for filling a liquid into holes that have aspect ratios of larger than 5:1 in a substrate which can be carried out by the operating steps of providing a substrate that is equipped with a plurality of through holes that have aspect ratios of larger than 5:1, each of the plurality of through holes has a first opening on a first surface of the substrate and a second opening on an opposite second surface of the substrate, contacting the first surface of the substrate with a vacuum plate equipped with a vacuum slot such that air is evacuated from the first opening of each of the plurality of through holes, and contacting the second surface of the substrate with a filler plate equipped with an injection slot such that a liquid can be injected into the second opening of each of the plurality of through holes which was evacuated by the vacuum plate.\nThe method for filling a liquid into holes that have aspect ratios of larger than 5:1 in a substrate may further include the step of contacting the first surface of the substrate with a vacuum plate and contacting the second surface of the substrate with a filler plate simultaneously. The method may further include the step of connecting the vacuum plate to the filler plate by a connection means which further includes an adjustment means of a veneer adjuster. The method may further include the step of adjusting the veneer adjuster such that the injection slot lags behind the vacuum slot when a substrate is slidingly displaced inbetween the filler plate and the vacuum plate. The method may further include the step of slidingly engaging a substrate between the filler plate and the vacuum plate, or slidingly moving a substrate inbetween the filler plate and the vacuum plate by a motor means, or the step of slidingly moving the substrate inbetween the filler plate and the vacuum plate at a speed between about 25 mm/min and about 250 mm/min.\nThe method for filling a liquid into high aspect ratio holes in a substrate may further include the step of covering substantially all the through holes in the substrate by the injection slots situated in the filler plate when the substrate is placed between the filler plate and the vacuum plate. The method may further include the step of connecting the vacuum plate to the filler plate in such a way that the vacuum plate is situated below the filler plate. The method may further include the step of providing a pressurized liquid reservoir in fluid communication with the injection slot, or the step of feeding a liquid to the injection slot in the filler plate from a liquid reservoir.\nThe method may further include the steps of feeding a molten solder in a liquid form to the injection slot in the filler plate, injecting the molten solder into the through holes, and solidifying the molten solder and forming via contacts. The method may further include the step of connecting a vacuum evacuation source to the vacuum slot in the vacuum plate such that air is evacuated from the first opening of each of the plurality of through holes.\nIn another preferred embodiment, the present invention provides an apparatus for filling an electrically conductive material into high aspect ratio holes in an electronic substrate which includes a first plate equipped with an injection slot for delivering an electrically conductive material, a second plate equipped with a vacuum slot for evacuating air, means for fixing the first and second plates together in a parallely face-to-face, spaced-apart relationship forming a dispensing unit for intimately engaging an electronic substrate therein, the electronic substrate may have holes with aspect ratios larger than 5:1 therein and the holes are in fluid communication with the injection slot and the vacuum slot when intimately engaged between the first and second plates, and drive means for causing sliding engagement and relative motion between the dispensing unit and an electronic substrate positioned therein.\nIn the apparatus for filling an electrically conductive material into high aspect ratio holes, the apparatus may further include an electrically conductive material reservoir for feeding the material into the injection slot. The apparatus may further include vacuum evacuation means for withdrawing air from the vacuum slot. The fixing means may further include adjustment means for adjusting a relative position of the first plate to the second plate, the adjustment means may be a veneer adjuster. The injection slot in the first plate may be adjusted behind the vacuum slot in the second plate, or the injection slot in the first plate may be adjusted to the same axial position of the vacuum slot in the second plate. The electrically conductive material delivered may be a molten solder at the lower end of the viscosity range, or a conductive polymer at the higher end of the viscosity range. The drive means drives the electronic substrate and cause it to slide between the first and second plates. The drive means may be a motor means.\nThe present invention is still further directed to a method for forming interconnects in a substrate for a display panel which can be carried out by the steps of first providing a substrate for a display panel that has a plurality of interconnect holes therein, each of the plurality of interconnect holes may have a first opening on a first surface of the substrate and a second opening on an opposite second surface of the substrate, withdrawing air from the plurality of interconnect holes by contacting the holes at the first opening with a vacuum slot provided on a vacuum plate, and filling the plurality of interconnect holes with a molten solder by contacting the holes at the second opening with an injection slot provided on a filler plate.\nIn the method for forming interconnects in a substrate for a display panel, the method may further include the step of solidifying the molten solder in the plurality of interconnect holes and forming the interconnects. The method may further include the step of contacting the plurality of interconnect holes with the vacuum slot and the injection slot simultaneously. The method may further include the step of forming the interconnects and placing a voltage potential on pixel display elements formed on the display panel. The method may further include the step of mounting the vacuum plate and the injection plate together in a parallely face-to-face, spaced-apart relationship with a cavity formed therein for slidingly engaging the substrate in abutting contact. The method may further include the step of adjusting a relative axial position of the vacuum plate to the injection plate, wherein the adjustment means may be a veneer adjuster. The method may further include the step of feeding a molten solder to the injection slot from a molten solder reservoir."}
{"text": "1. Field of the Invention\nThis invention relates, in general, to pictorial communications and, more specifically, to signal processing techniques for use with the electronic copiers, printers, and like devices.\n2. Description of the Prior Art\nDigital halftoning can be used in devices or systems, such as copiers and printers, where the pixels of the image data are represented by digital values. Halftoning is frequently used when the rendering device, or print engine, does not have the same gray-level capability as the input information. In halftoning, halftone cells are created wherein the overall density of the halftone cell is expected to be equal to the density of a group of pixels in the original input data. This provides the appearance of the same gray-level density for the area without the need for a gray-level rendering device. The most frequent use of digital halftoning is associated with reproducing gray-level input data with a binary output device. Although the overall density of the cell appears to be near the desired density, some sharpness information is lost from the original pixel data with many conventional halftone processing systems.\nError diffusion is another technique which may be used to convert gray-level values to binary values, or to gray-level values with a smaller number of usable levels. However, if error diffusion is used alone, the process stability of a single pixel dot at high resolution is a usual concern. In some processing, error diffusion and digital halftoning can be used together to provide the proper data for rendering the image on a binary device. This can provide good process stability and more overall-area gray-levels. If error diffusion is applied to each pixel of the halftone cell, much of the sharpness of the original data is preserved. However, pixel-to-pixel error diffusion is computationally intensive and can present major speed and hardware restraints. Error diffusion between halftone cells is much more economical as far as processing time is concerned, but a degree of sharpness is sacrificed with this arrangement.\nU.S. Pat. No. 4,680,645, issued on July 14, 1987, discloses a method for rendering gray-level image data by a binary device. According to this patent, pixel-to-pixel error diffusion is used in conjunction with a variable sized dot rendering system. The variable sizes are modulated by the error in the gray-levels diffused or propagated to other pixels. Halftone cells are not used in this system. The background section of the patent includes a general discussion on error diffusion at column 1, lines 42-64.\nIn order to obtain an improved digital halftoning system, it is desirable, and an object of this invention, to provide a system for converting gray-scale data into binary data which accurately represents the density and sharpness information in the original image with reduced computations."}
{"text": "1. Field of Invention\nThe present invention relates generally to catheter systems used in diagnosis and treatment of various body tissues and, more specifically, to steering systems incorporated into catheter systems for ablating cardiac tissue in the treatment of electrophysiological diseases.\n2. Description of the Related Art\nAs is well know, catheters provide medical professionals access to various interior regions of the human body, in a minimally invasive manner, in support of diagnosis and treatment. Catheters allow such professionals to place one or more medical instruments, pharmacological agents or other matter at a target tissue site. For example, in cardiac procedures in support of diagnosis and treatment of atrial fibrillation, catheters provide access to various chambers of the heart, carrying ablation devices which translate therein to such sites for ablation of specific cardiac tissue associated with atrial fibrillation.\nAblation of tissue, cardiac tissue for example, is directly related to the orientation of the ablation element from which energy sufficient to ablate biological tissue is emitted. For such procedures, precise control of the ablation device is desirable to ensure proper placement of the ablation element utilized in creation of one or more desired lesions. As a surgeon, or other medical professional, manipulates the proximal end of the catheter system, the distal end of the catheter must be responsive to such movement in a very predetermined, smooth-flowing and proportional way.\nProper placement of an ablation device is exasperated by the fact that some ablative energy technologies require energy transmission conduits which are bulky, or otherwise constructed from materials less flexible, making the distal portion of the catheter difficult to properly position. For example, distal portions of optical fiber or microwave based ablation systems, or catheter systems comprising an endoscopic device, may be more difficult to maneuver due to the lack of {flexibility in the transmission mediums utilized therein. As should be readily apparent, when the distal portion of an ablation catheter system is not properly positioned, ablative energy is not properly directed and applied to the target tissue, resulting in poor lesion formation. It is therefore essential that the ablative device be able to be manipulated and sufficiently controlled to be properly positioned to transfer the requisite energy to ablate biological tissue and create a desired lesion therein."}
{"text": "Analog-to-digital converters (ADCs) convert time-discrete analog input values to a digital form. A type of ADC, the successive approximation register (SAR) ADC, digitizes the analog input values using a successive approximation search algorithm. While the internal circuitry of the SAR ADC may run at a higher frequency (such as several megahertz (MHz), for example), the sample rate of the SAR ADC is generally a fraction of that frequency (such as several kilohertz (kHz), for example) due to the successive approximation search algorithm used. For example, normally each bit of the SAR ADC is fully realized prior to proceeding on to the next bit.\nSince modern ADC devices are highly precise, testing them calls for even more precise stimuli for the results to be of relevance. Providing these high precision stimuli can be expensive. Additionally, in order to reach the required high precision, the stimuli generators often need long settling times, resulting in longer test times.\nOften, front-end and back-end testing using automatic test equipment (ATE) is performed as part of ADC production. For production tests using ATE, the ADCs can be tested with architecture-independent methods, such as the linear ramp test or the histogram test. With the linear ramp test method, a slowly changing high-resolution high-linearity voltage ramp is applied to the input of the ADC, and the conversion results are used to determine the ADC transfer curve. The histogram test method uses a predetermined input signal as input to the ADC, and statistically deduces the transfer curve from the observed code histogram.\nBoth methods often use a minimum number of hits per code to provide usable results. Consequently, the runtime for an n-bit binary ADC grows linearly with the number of possible output codes, which is 2^n. For high resolution ADCs, this leads to long test times and, as a consequence, high test costs.\nA compromise for easier stimulus generation may be achieved by using a first-order resistor-capacitor low pass filter. The inherent mismatch of the RC time constant is measured first, and then the transfer curve can be obtained. As the charge or discharge speed of the capacitance must be adapted to the fastest changing part of the voltage curve, the measurement takes longer than the linear ramp test. Additionally, the high time constant can require an unfeasibly large (external) capacitance.\nExpensive stimuli generators and long settling times can lead to increased test costs. Further, when tests are lengthy and expensive, the user may have no simple means to check that the circuits are fully functional during field operation, including throughout the product lifetime. For example, some methods are used for production testing, and are prohibitive for use in the field. This can lead to a violation of specified functional safety requirements of some automotive applications, for example.\nOther proposed methods often result in substantial chip area overhead and can also increase the BOM for the customer."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of display, and in particular to a thin film transistor (TFT) backplane and manufacturing method thereof.\n2. The Related Arts\nIn the field of display technology, liquid crystal display (LCD) and organic light emitting diode (OLED) display and other flat panel display technology have gradually replaced CRT monitors, wherein the OLED provides the advantages of self-luminous, low driving voltage, high luminous efficiency, short response time, high clarity and contrast, wide viewing angle of near-180°, wide operating temperature range, enabling flexible display and large full-color display, and is recognized as the technology with most potential.\nThe driving types of OLED can be divided according to the passive matrix OLED (PMOLED) and active matrix OLED (AMOLED). The low temperature poly silicon (LTPS) thin film transistor (TFT) has gained attention from the industry in high resolution AMOLED technology. Compared with the amorphous silicon (a-Si), LTPS TFT has higher carrier mobility, and the device has fast response and good stability to meet the requirements of high resolution AMOLED display.\nThe gate insulation (GI) commonly used in known LTPS TFT backplane of AMOLED usually uses the double-layered structure of silicon dioxide/silicon nitride (SiO2/SiNx), wherein the SiO2 layer contacts the polysilicon active layer, and the SiNx layer contacts the gate. Compared to SiO2 layer, SiNx layer provides better barrier for mobile ions, such as, sodium ions (Na+), potassium ions (K+), has a greater dielectric constant, can be made thinner than the SiO2 layer under the same insulating capacity, has a higher hydrogen (H) content, and has a passivation effect on dangling bond in polysilicon. However, the SiNx layer in contact with the polysilicon active layer has a poor interface due to the stress. Hence, a SiO2 layer is often deposited plus a layer of SiNx layer to form the gate insulating layer.\nThe thicker the SiNx layer is, the better the capabilities are as barrier to mobile ions and passivation effect. However, the reliability of the TFT elements will decrease. This is because the gate will continually inject the carrier stream into the SiNx layer, which will destroy the SiNx layer, so that the quality of the SiNx layer is deteriorated, resulting in reduced reliability of the TFT element. Therefore, under normal circumstances, SiNx layer is not thick. In addition, during fabricating the gate in the etching process, the SiNx layer is often over-etched. If a protective layer is superimposed over the SiNx layer simply to protect the SiNx layer, the thickness of the gate insulating layer corresponding to storage capacitor area is increased, resulting in performance degradation in capacitive storage. The only solution is to increase the capacitance area and sacrifice the opening ratio to ensure the capacitance storage performance."}
{"text": "1. Technical Field\nThe embodiments herein generally relate to emergency control mechanisms, and, more particularly, to an emergency safety shutoff system for a glovebox air intake.\n2. Description of the Related Art\nEnvironmental accidents, toxic chemical spills, and testing to support verification inspections are examples of situations where chemical and biological testing may be required. Gloveboxes are used for such testing. The availability of safe, reliable, and efficient containment vessels, such as gloveboxes and fume hoods, for conducting environmental tests help to protect the investigators and preserve the evidentiary value of the collected materials and the on-site tests.\nThe glovebox serves as engineering control and primary protection during a hazardous test operation. During the test, in the event of pressure loss due to power failure or exhaust blower malfunction, there is a potential for containment loss through the air intake valve which would allow for hazardous materials to enter into the operators breathing zone. It is desirable to eliminate the potential hazard with an improved glovebox system."}
{"text": "In U.S. Pat. No. 6,888,196, a conventional structure of power semiconductor device is disclosed, as shown in FIG. 1, wherein an N-channel trench MOSFET comprising a plurality of trenched gates 110 surrounded by n+ source regions 112 encompassed in P body regions 114 is formed in an N epitaxial layer 102 over an N+ substrate 100. To connect said source regions 112 and said body regions 114 to a source metal 122, a trenched source-body contact 118 with vertical sidewall is employed penetrating through a contact interlayer 120, said n+ source regions 112 and extending into said P body regions 114. Furthermore, a p+ body ohmic contact doped region 116 is implanted surrounding bottom of said trenched source-body contact to decrease a contact resistance between said P body regions 114 and said trenched source-body contact 118.\nThe conventional structure in FIG. 1 is accoutering a technical difficulty which is that avalanche always occurs near bottom of said trenched gates 110, causing a hazardous condition to the power semiconductor device. As we all know that, in the trench MOSFET shown in FIG. 1, a avalanche current Iav (illustrated in FIG. 1) flows between said trenched gates 110 and said source-body contact 118, triggering turning-on of a parasitic bipolar transistor (illustrated in FIG. 1) when Iav*Rb>0.7V, wherein Rb is a resistance between said p+ body ohmic contact doped region 116 and channel region near said trenched gates 110. As is known to all that, the doping concentration of said p+ body ohmic contact doped region 116 is higher than that of said P body region 114 (please refer to FIG. 2 for Y1-Y1′ cross section of FIG. 1), which is helpful to decrease resistance Rb, however, as the sidewall of said trenched source-body contact is perpendicular to the front surface of said N epitaxial layer 102, when carrying out implantation through a contact trench, said p+ body ohmic contact doped region 116 can be formed only surrounding bottom of said trenched source-body contact, resulting in a high resistance Rb underneath said n+ source regions 112. Therefore, said parasitic bipolar transistor is easily to be triggered turning on due to the high resistance Rb, thus weakening the avalanche capability of the trench MOSFET.\nFIG. 3 shows another trench MOSFET in prior art disclosed in U.S. Patent No. 20080890357. Comparing to FIG. 1, the trench MOSFET in FIG. 3 comprises a plurality of trenched gates 130 having terrace gate structure for gate resistance reduction, wherein top surface of gate conductive layer filled in gate trenches is higher than the sidewall. However, the limitation of poor avalanche capability discussed above is still pronounced in this structure due to the easily turning-on of a parasitic bipolar transistor and the occurring of avalanche near bottom of said trenched gates 130.\nFor other power semiconductor power device, for example trench IGBTs, the same disadvantage of poor avalanche capability is also affecting the performance of the power semiconductor device.\nAccordingly, it would be desirable to provide new and improved power semiconductor devices to avoid the constraint discussed above."}
{"text": "This invention is in the field of terminally guided, anit-armor missiles. Heretofore known terminal guidance missile systems have used proportional navigation with limited trajectory shaping for high accuracy against moving targets. This limited use of trajectory shaping results in either a flat approach trajectory which has reduced warhead penetration or a lofted or ballistic like trajectory. The ballistic like trajectory is often unable to perform well when low cloud cover condition exists. The including of a gimbal angle regulator and boost/cruise trajectory shaping allows low altitude ground or air launch, climb to cruise altitude under a low cloud cover, and then dive onto the target thereby achieving a high probability of penetrating the armor of the target. This invention is not limited to tracker systems that must acquire the target prior to launch such as an infrared imaging seeker, but works equally well with systems that can acquire the target after launch such as laser semiactive systems. The use of conventional known guidance schemes cannot accomplish the high probability of accurate hit concurrent with control of the dive angle for maximum warhead performance. Any ballistic like trajectory for terminally guided missiles must reacquire after descending through the cloud cover adversely affecting the probability of hit and limiting the controllability of the impact attitude."}
{"text": "The present invention relates to fittings for securing articles in a desired position and, more particularly, to a tie-down connector fitting for securing automobiles and the like to the bed of a transporter, such as a railroad car or truck.\nAutomobiles and other vehicles are normally transported from the factory to a location near the point of purchase by transporters, including trucks and railroad cars. In the past, various means of securing vehicles to such transporters have been used, the most popular method being a tie-down assembly attaching the undercarriage of the vehicle to the bed of the transporter. Generally, prior tie-down assemblies consist of a hook section and a single chain extending therefrom to an idler assembly movable along a track formed in the floor or bed of the transporter to accommodate vehicles of different lengths. In securing a vehicle, the idler assembly is moved along the track to a position beneath the vehicle, where the hook is attached to an opening or slot in the undercarriage of the vehicle. The chain is attached at its free end to a conventional ratchet, movable within the track, which tightens the chain for securing the vehicle to the bed of the transporter for shipment. This procedure was often repeated at each of the four corners of the vehicle as a safety precaution to avoid potentially dangerous and damaging movement of the vehicles should one of the chains break.\nIn an effort to reduce material and labor costs in tying down vehicles for shipment, shippers have begun to use single tie-down assemblies at the front and rear of the vehicles. Such an arrangement, however, introduces the danger of a loose vehicle on the transporter should one of the chains break. Accordingly, improved vehicle tie-down systems have been developed to safeguard against the possibility of chain failure, such as that disclosed in U.S. Pat. No. 4,022,134 to Krokos, issued May 10, 1977, which shows multiple chains attached to a clevis 41, pivotally pinned to the center of a connector fitting formed in various configurations. A further improvement is shown in patent application Ser. No. 703,594, filed July 8, 1976, now U.S. Pat. No. 4,044,691, issued Aug. 30, 1977, by Donald C. Kilgus and entitled \"Vehicle Tie-Down System for Rail Cars\".\nThese improved tie-down systems generally provide a connector fitting which attaches to a pair of chains or at a point along the length of a single chain in a triangular configuration. In such systems, each separate chain or the ends of a single chain are secured to the transporter so that one will hold the vehicle alone, should the other break. It has been discovered in using such tie-down systems, however, that a great deal of stress and strain is placed on the fitting itself, and particularly on the end link of the chain which attaches to the fitting. The tensile and torsion stresses applied to the chains and fitting by the motion of the vehicle on the transporter, especially in starting and stopping, have caused failures in both after relatively short periods of use.\nAccordingly, it is desirable to provide a tie-down system wherein the chains may be separated from the fitting should one fail before the other. However, in many prior art designs, such as that disclosed in U.S. Pat. No. 4,049,229 to Peisner, issued Sept. 20, 1977, the chains are integral with the fitting and may not be separated therefrom. In such designs, if either one of the fitting or chains prematurely fail, the entire connector fitting, including the chains, must be replaced."}
{"text": "Clamping claws are generally used to secure an object in a recess of a support. The clamping claws have a securing leg and a clamping leg, with the clamping claw in the installed state being secured on the support by way of its securing leg by means of a bolt and clamping the object in the recess by means of its clamping leg.\nSuch clamping claws are used, for example, as part of a cutting tool for machining metallic materials, with the support also being referred to as a gripping holder. The object that is clamped by means of the clamping claw in this connection is a cutting plate made from a ceramic material, hard metal or other material that is suitable for machining.\nThe underlying aim of the invention is to convert the application of axial force applied to an object by a clamping claw in part into an application of force that acts transversely thereto so that the object is pressed not only onto the support, but also onto the side wall of the recess.\nThe clamping claw includes two portions, namely a clamping base and a clamping finger.\nThe clamping base sits on the support in the installed state of the clamping claw.\nThe clamping finger has a securing leg and a clamping leg, and in the installed state the securing leg sits on the clamping base over an oblique plane formed as a sliding surface and clamps the object with its clamping leg.\nDuring clamping the bolt presses the securing leg of the clamping finger onto the clamping base and secures both on the support, in which case as a result of the oblique plane that is formed as a sliding surface the securing leg and thus the clamping leg of the clamping finger execute a lateral movement transversely in relation to the longitudinal axis of the clamping base.\nThe application of force applied to the object by the clamping claw is split into an axial component and a component that acts transversely thereto.\nSo that the component that acts transversely in relation to the axial component draws the object into the recess of the support, a trough, for example, into which, for example, a cam on the clamping leg of the clamping finger engages, is arranged in the object.\nThe clamping base is advantageously secured on the support in a secure manner with respect to rotation, but so that it can be displaced in the clamping direction. To this end, at least one projection that extends parallel to the longitudinal axis of the clamping base and engages into a longitudinal slot in the support that is matched thereto is arranged on the clamping base.\nIn a preferred embodiment, the clamping base at its end that faces the securing leg of the clamping finger has an oblique plane that is formed as a sliding surface, and the securing leg of the clamping finger at its end that faces the clamping base likewise has an oblique plane that is formed as a sliding surface. Both sliding surfaces are matched to each other and are formed so that they slide one on top of the other.\nThe clamping base in the installed state preferably rests with a first guide surface against a second guide surface on the securing leg so that the securing leg of the clamping finger is arranged on the clamping base in a secure manner with respect to rotation. Since the clamping base also sits on the support in the installed state in a secure manner with respect to rotation, the clamping finger is also arranged on the support in a secure manner with respect to rotation. This is necessary in order to achieve correct clamping of the object.\nAdvantageously, the clamping base has a projecting flange with which the clamping base sits on a recess in the support in the installed state.\nThe clamping claw is preferably a part of a cutting tool for machining metallic materials, in which case the support is a gripping holder and the object is a cutting plate.\nIn an advantageous embodiment, arranged on the clamping leg of the clamping finger there is a raised portion, like a cam for example, which in the installed state engages into a corresponding trough in the object or in the cutting plate.\nIn a preferred development in accordance with the invention, the raised portion or the cam on the clamping leg of the clamping finger is formed as a circular ring which in the installed state engages into an object or a cutting plate that has a spherical or circular raised portion and in so doing embraces the raised portion.\nThe invention is explained in greater detail in the following with the aid of six figures."}
{"text": "This invention is concerned with an adapter to enable a dry cell battery to be used in an application in which it might be charging or discharging for which it is electrically acceptable but for which it is physically too small.\nBatteries in flash lights, portable radios or so called chargers and other equipment are usually supported in a tube or clips or other mechanical contrivances so that they are securely held in electrical contact with the terminals of the circuit or element which they are to supply or from which they are to be supplied (charged). The support means may hold a single battery between the terminals or a plurality of batteries in end to end relationhip may be so held but almost always the support means are arranged to receive a battery of particular transverse cross sectional dimensions.\nAs is well known, there are generally available several differently dimensioned batteries which have similar electrical discharge and charging characteristics. It frequently occurs than when a physically large battery is needed there is available only a battery of a smaller size but one which is electrically acceptable.\nFurther, as batteries are reduced in size with advancing technology, it becomes technically and economically desirable to use the new and smaller batteries in equipment designed specifically to receive older and larger batteries. For example, a nickel/cadmium AA battery has similar electrical discharge characteristics to a zinc/carbon D battery but of course as it is rechargeable it is desirable to be able to use such a battery in place of the larger zinc/carbon battery.\nMajor differences between existing large and small batteries is their working lives. Many pieces of equipment such as, for example, flash lights are designed to receive large batteries that fit them for long and continuous use but in practice those pieces of equipment may be used for only short periods with long intervals of non-use between those periods. During those periods of non-use the battery deteriorates or, because of faulty terminals in the equipment, discharges. In many cases these long periods of non-use can be foreseen but because of the physical design characteristics of the pieces of equipment one is obliged to use a large and expensive battery knowing that were its use physically possible a small, cheaper battery having a shorter working like would be quite adequate."}
{"text": "Fuel injected engines employ fuel injectors, each of which delivers a metered quantity of fuel to an associated engine cylinder during each engine cycle. Prior fuel injectors were of the mechanically or hydraulically actuated type with either mechanical or hydraulic control of fuel delivery. More recently, electronically controlled fuel injectors have been developed. In the case of an electronic unit injector, fuel is supplied to the injector by a transfer pump. The injector includes a plunger which is movable by a cam-driven rocker arm to compress the fuel delivered by the transfer pump to a high pressure. An electrically operated mechanism either carried outside the injector body or disposed within the injector proper is then actuated to cause fuel delivery fuel to the associated engine cylinder.\nIn prior fuel injector designs, high pressure fuel is conducted through passages which are located outside of a central recess containing a solenoid which operates a valving mechanism. The passages are located close to the outer surface of the fuel injector and are formed by drilling intersecting holes. After drilling, portions of some of the holes must be filled with plugs. These passages and plugs are subjected to very high fluid pressures, thus requiring careful design, thus increasing complexity and cost.\nIn addition to the foregoing, because the high pressure passages are located outside of the solenoid, the size of the solenoid is necessarily limited, thereby limiting the available solenoid force.\nStill further, a prior type of fuel injector utilizes a direct operated check valve, which includes upper and lower valve seats which must be precisely aligned for proper operation. Manufacturing and assembly tolerances must, therefore, be kept tight, further increasing cost."}
{"text": "Physicians have often used endoscopes to examine, to biopsy, and to ablate the tissue of patients within lumens such as the esophageous and the bowel. These procedures include esophageal duodenoscopy, (EGD), colonoscopy, and polypectomy. Endoscopes used in these procedures can be several feet long and generally comprise one or more instrument channels and optical fiber bundles. The instrument channels and optical fiber bundles open into the body at the distal end of the endoscope and are generally parallel to the axis of the flexible endoscope. Physicians place flexible instruments through the instrument channels while visualizing and illuminating a site using the optical fiber bundles. The instruments have end-effectors at the distal end for performing useful therapeutic work to tissue. The instruments also provide at their proximal ends operating mechanisms for actuating the end-effectors. Such instruments placed through an instrument channel of an endoscope may include biopsy forceps for tissue sampling, electrical wires for radiofrequency ablation, or tubes used for irrigation, gas transfer, particular matter transfer, and suction.\nA physician performing a therapeutic procedure with the use of an endoscope places a long, flexible instrument through the endoscope's instrument channel and then positions the instrument near the site within the body lumen where a therapeutic procedure is to be performed. The physician grasps the endoscope with one hand and introduces the flexible instrument from an entrance at the proximal end of the endoscope with the other. An assistant usually holds the proximal end of the flexible instrument and operates the mechanism at the proximal end of the instrument to actuate the end-effector.\nLocating the end-effector within the lumen presents difficulties to the physician. The flexible instrument emerges from the instrument channel of the endoscope in a direction parallel to the axis of the endoscope. Many times work needs to be performed by the instrument at the inner wall of a body lumen, and at an angle to the axis of the endoscope. Much tedious maneuvering of both endoscope and instrument is needed to place the end-effector into a position to perform useful work. The flexible nature of the instrument creates difficulties in locating the end-effector at the needed position on the inner wall of the body lumen. Rotating the endoscope to accommodate the instrument causes the image on the monitor to rotate creating visualization difficulties for the physician.\nWhen the task to be performed is a biopsy, a tissue sample or specimen must be transported away from the work site for collection. To collect the specimen, the physician pulls the entire instrument from the endoscope, removes the specimen, places the specimen into a collection jar, and replaces the instrument into the endoscope. The physician must then manipulate the instrument into a new position to take another specimen."}
{"text": "A dicing apparatus, which performs a cutting or groove-forming process on a workpiece such as a wafer on which semiconductor elements or electronic components are formed, includes at least: a rotary blade which is formed of a thin grindstone called a blade and is rotated by a spindle at high speed; a worktable which holds the workpiece; and X, Y, Z, and θ moving shafts which change a relative position between the worktable and the blade. At the time of processing the workpiece, a cutting fluid for cooling or lubrication is supplied from a nozzle to the rotating blade or a processing point at which the workpiece and the blade contact each other, and the cutting or groove-forming process is performed on the workpiece by the operations of the respective moving shafts.\nFIG. 8 illustrates a conventional example of the dicing apparatus. A dicing apparatus 70 includes a processing part 75. The processing part has: high-frequency motor built-in type spindles 72, 72 respectively equipped with a blade 71 and a wheel cover (not shown) on their tips, and disposed so as to be opposed to each other; a worktable 73 which adsorbs and holds a workpiece W; and an image pickup device 74 formed of a microscope, a CCD camera, or the like for picking up an image of the workpiece W. In addition to these components, the dicing apparatus 70 includes: a cleaning part 76 which spin-cleans the processed workpiece W which has been processed by the processing part 75; a load port 77 onto which a cassette that houses a large number of the workpieces W mounted on a frame F is set; a conveyance device 78 which conveys the workpiece W; a controller 79 which controls the operations of the respective parts; and the like.\nAs illustrated in FIG. 9, the structure of the processing part 75 includes an X table 83 which is guided by X guides 81, 81 provided to an X base 80 and is driven by a linear motor 82 in an X direction indicated by X-X of the figure. A worktable 85 is provided to the X table 83 via a rotary table 84 which rotates in a θ direction.\nOn the other hand, Y tables 88, 88 which are guided by Y guides 87, 87 and are driven by a stepper motor (not shown) and a ball screw (not shown) in a Y direction indicated by Y-Y of the figure are provided on a side surface of a Y base 86. A Z table 89 which is driven by a drive device (not shown) in a Z direction indicated by Z-Z of the figure is provided to each of the Y tables 88. The high-frequency motor built-in type spindle 72 which has the blade 71 on its tip is fixed to the Z table 89. The processing part 75 has the structure as described above, and hence the blade 71 is step-fed (fed in a stepwise manner) in the Y direction and also is fed while cutting-in (cutting-in feed) in the Z direction, and the worktable 73 is fed while cutting (cutting feed) in the X direction.\nThe spindles 72 are both rotated at a high speed of 1,000 rpm to 80,000 rpm, and a supply nozzle (not shown) which supplies the cutting fluid into which the workpiece W is to be immersed is provided in the vicinity of the spindles 72.\nAn electrodeposition blade obtained by electrodepositing diamond abrasive grains or CBN abrasive grains with nickel, a metal-resin bonding blade obtained by bonding with a resin mixed with a metal power, or the like is used as the blade 71. The dimensions of the blade 71 are variously selected depending on the processing type. In a case where a general semiconductor wafer is diced into the workpiece, a blade having a diameter of approximately 50 mm and a thickness of approximately 30 μm is used.\nIn addition, the controller which controls the operations of the respective parts of the dicing apparatus 70 includes a CPU, a memory, an input/output circuit part, various control circuit parts, and the like, and is incorporated inside a pedestal of the dicing apparatus 70. As the dicing apparatus having the above-mentioned structure, for example, a dicing apparatus disclosed in Patent Document 1 has been proposed.\nIn the dicing apparatus 70 having the above-mentioned structure, as a stage prior to the processing, an alignment operation is performed, in which an image of the workpiece W is picked up by the image pickup device 74 and the position of the workpiece W is aligned with the blade 71. Also during the processing, if necessary, an image of the workpiece W is picked up by the image pickup device 74 as appropriate to check a processing condition. However, in the dicing apparatus 70, only one worktable 73 onto which the workpiece W is set and only one image pickup device 74 which picks up an image of the workpiece are provided. Therefore, during the processing of one workpiece W, it is not possible to set another new workpiece W onto the worktable 73 to perform the alignment by the image pickup device 74. Accordingly, the utilization rate of the image pickup device is low, and in addition, the utilization rate of the entire dicing apparatus is also deteriorated.\nAs a solution to such a problem, Patent Document 2 discloses a dicing apparatus in which two blades, two worktables, and two image pickup devices are provided. In this dicing apparatus, a workpiece which is set onto one worktable is aligned by one image pickup device, while a workpiece which is set onto another worktable is aligned by another image pickup device. Then, the workpiece which is set onto the one worktable is processed by one blade, while the workpiece which is set onto the another worktable is processed by another blade. Alternatively, the workpiece which is set onto the one worktable or the another worktable is processed by both the one blade and the another blade. In this manner, two workpieces are aligned individually by the two image pickup devices, and one workpiece is processed by the two blades, whereby the utilization rate of the dicing apparatus is improved.    Patent Document 1: Japanese Patent Application Laid-Open No. 2002-280328    Patent Document 2: Japanese Patent Application Laid-Open No. 2006-156809"}
{"text": "It is well known that photovoltaic cells use specially prepared semiconductor junctions to convert energy from sunlight to electricity.\nIt is also known that a specific photovoltaic semiconductor junction utilizes, for conversion to electricity, only a portion of the spectrum of energy available in the available light. For example, the conversion of sunlight to electrical energy using well known silicon photovoltaic cells is strongly dependent upon the conversion of light with energy at or above 1.1 electron volts while the lower energy light also present is, at least in part, converted to heat instead of electricity. The heat generated can reduce the efficiency of the silicon cell for the conversion of the higher energy light to electricity.\nSome of the sunlight which penetrates the solid light transmitting member of a conventional photovoltaic array is lost such as by passing through an interstice between adjacent photovoltaic cells, reflection back out of the member, absorption, and the like.\nAccordingly, current photovoltaic arrays receive more energy input from incident light than they retain for conversion into electrical output, and it is highly desirable to increase the amount of light an array retains for such conversion.\nOne known approach for increasing the capacity of photovoltaic arrays to convert light energy to electricity is to employ one or more solid luminescent agents in the solid light transmitting member. Such agents, when exposed to sunlight, take in light from one direction and emit lower energy light in numerous directions. Examples of such agents are organic dyes such as the dyes heretofore employed in scintillation counters, lasers, and the like.\nThe particular luminescent agent or agents employed in conjunction with an array of specific photovoltaic cells are chosen, inter alia, for their ability to emit light in an energy level range which suits the conversion characteristics of that cell. This way, a portion of the light that would otherwise be lost for electrical generation by transmission, reflection, and the like is transformed by the luminescent agent into multi-directional light that is more readily retained in the light transmitting member and which the photovoltaic cell or cells can readily convert into electricity, thereby increasing the overall efficiency of the array.\nA photovoltaic array which employs this approach is referred to as a luminescent photovoltaic array. Such an array usually employs fluorescent dyes, fluorescence being that species of luminescence wherein the emitted light is usually in the visible spectrum. However, other species of luminescence exist. Phosphoresence (light emission continues after the stimulating light has stopped) is one such species.\nIt should be understood that this invention covers all species of luminescence, as well as all types of luminescent agents.\nHeretofore the prior art has required the photovoltaic cell or cells employed to be mounted on the thin edge of a solid, plastic or glass, light transmitting luminescent member. Such a device is fully and completely disclosed in Applied Optics, Volume 15, No. 10, Pages 2299 and 2300, October, 1976, the disclosure of which is incorporated herein by reference."}
{"text": "1. Technical Field\nIn general, the disclosed embodiments relate to devices and methods for the introduction and subsequent evaluation of therapeutic agents to biological tissue, and in particular to the simultaneous introduction of a plurality of agents to the tissue in vivo.\n2. Description of the Related Art\nNumerous cancer-related therapeutics are under phase I or phase II clinical trial and evaluations at any particular time; however, most of them will fail to advance. In fact, it is estimated that more than 90% of cancer-related therapeutics will fail phase I or II clinical trial evaluation. The failure rate in phase III trials is almost 50%, and the cost of new drug development from discovery through phase III trials is between $0.8 billion and $1.7 billion and can take between eight and ten years.\nIn addition, many patients fail to respond even to standard drugs that have been shown to be efficacious. For reasons that are not currently well understood or easily evaluated, individual patients may not respond to standard drug therapy. One significant challenge in the field of oncology is to exclude drug selection for individual patients having cell autonomous resistance to a candidate drug to reduce the risk of unnecessary side effects. A related problem is that excessive systemic concentrations are required for many oncology drug candidates in efforts to achieve a desired concentration at a tumor site, an issue compounded by poor drug penetration in many under-vascularized tumors (Tunggal et al., 1999 Clin. Canc. Res. 5:1583).\nClearly there is a need in the art for improved devices and methods for testing and delivering cancer therapies, including improved methodologies for performing efficient pre-clinical and clinical studies of candidate oncology medicines, and for identifying therapeutics having increased likelihood of benefiting individual subjects. The present invention addresses these and similar needs, and offers other related advantages."}
{"text": "The senile plaque and congophilic angiopathy are abnormal extracellular structures found in abundance in brain of patients with Alzheimer's disease. The biochemical composition of these structures has been extensively studied to better understand their possible role in the pathogenesis of this dementing disease. The mature senile plaque is a complex structure, consisting of a central core of amyloid fibrils surrounded by dystrophic neurites, axonal terminals and dendrites, microglia and fibrous astrocytes. See D. Selkoe Neuron 6, 487-498 (1991). The amyloid core of the senile plaque surrounding blood vessels, producing the congophilic angiopathy, is a peptide of 39 to 43 amino acids termed the .beta.-Amyloid (A.beta.) peptide. G. Glenner and C. Wong, Biochem. Bioshys. Res. Comm. 120, 885-890 (1984). A.beta. peptide is found in brain in Alzheimer's disease, Down's syndrome, hereditary cerebral hemorrhage of the Dutch type, and in old age. K. Kosik, Science 256, 780-783 (1992). A.beta. is produced by abnormal proteolytic processing of a larger protein, the amyloid precursor protein (APP). See K. Bayreuther and C. Masters, Brain Path. 1, 241-251 (1991).\nThe senile plaque and congophilic angiopathy contain proteins in addition to .beta.A peptide. APP itself, among others, has been identified in the senile plaque by histochemical studies employing antibodies recognizing either the amino- and carboxy- termini of the precursor protein. See, e.g., F. Tagliavine et al., Neurosci. Lett. 128, 117-120 (1991); C. Joachim et al., Amer. Jour. Path. 138, 373-384 (1991); The mechanisms by which these proteins aggregate in the extracellular space to associate with the senile plaque and congophilic angiopathy are not known.\nApolipoprotein E (ApoE) performs various functions as a protein constituent of plasma lipoproteins, including its role in cholesterol metabolism. It was first identified as a constituent of liver-synthesized very low density lipoproteins which function in the transport of triglycerides from the liver to peripheral tissues. There are three major isoforms of ApoE, referred to as ApoE2, ApoE3 and ApoE4 which are products of three alleles at a single gene locus. Three homozygous phenotypes (Apo-E2/2, E3/3, and E4/4) and three heterozygous phenotypes (ApoE3/2, E4/3 and E4/2) arise from the expression of any two of the three alleles. The most common phenotype is ApoE3/3 and the most common allele is E3. See Mahley, R. W., Science 240:622-630 (1988).\nThe amino acid sequences of the three types differ only slightly. ApoE4 differs from ApoE3 in that in ApoE4 arginine is substituted for the normally occurring cysteine at amino acid residue 112. The most common form of ApoE2 differs from ApoE3 at residue 158, where cysteine is substituted for the normally occurring arginine. See Mahley, Science, supra.\nWhile there has been considerable research into the mechanisms underlying Alzheimer's disease, there continues to be an ongoing need for new ways to investigate and combat this disorder."}
{"text": "1. Field of the Invention\nThe present invention relates to a mounting arrangement for fixedly mounting one end of a piezoelectric element which comprises a carrier strip having at least one ceramic strip applied thereto and which element operates in a bending mode about the point of mounting for producing essentially linear movement for a part of a moving component operatively connected to the other end of the element.\n2. Prior Art\nPiezoelectric elements or piezoelectric benders which operate in a bending mode are well known in prior art. For example, an element having a metal carrier strip disposed between a pair of strips of piezoelectric material, such as lead-zirconium titanate, with electrodes for applying electrical current to the ceramic strip have been suggested. One proposed use of such elements is by rigidly mounting one end of the element in a body and by attaching the other end to a part of a movable component so that during a bending mode, the element creates substantially linear movement of the part attached to the other or free end. Various arrangements have been proposed to form so-called piezoelectric motors which have one or more of the elements arranged in parallel and mounted in the above manner. It has also been suggested to vary either the particular orientation of the ceramic layer or the application of the electrical current so that the element will bend in a substantial S configuration during operation.\nWhile the elements used either alone or in a motor arrangement have proven very useful, they have suffered failure due to the high stress applied to the element adjacent to the rigid mounting point at the one end or in the case of an S bending operation at both ends. In the case of the element which is rigidly mounted at each end and has an S shape bending mode, failure often occurs at either of the fixed or rigid ends."}
{"text": "Toothpaste dispensers are well known in the art. Most existing toothpaste dispensers are in the form of simple tubes that dispense toothpaste when users manually squeeze the tubes. Some of these tubes are formed of flexible plastic, and older ones of these tubes were formed of a soft, malleable metal. Alternatively, some toothpaste dispensers exist in a pump form. Simpler versions of these pump dispensers are manual and a user operates them by depressing the top of the dispenser, and other variants utilize a pump lever on the dispenser. Such toothpaste dispensers are not user-friendly for people who have arthritis in their hands or have other physical limitations such as Parkinson's disease.\nIn addition, such traditional tube and pump toothpaste dispensers are messy. It is easy to squeeze a tube of toothpaste or to press a manual pump type toothpaste dispenser with an incorrect amount of pressure and eject too much toothpaste that falls on a counter or, worse, on the floor. Further, the caps of tube type toothpaste dispensers often become messy and covered with toothpaste that dries out and thereby makes a worse mess.\nIn view of the shortcomings described above automatic toothpaste dispenser that can dispense toothpaste without manual operation have come into vogue. They are easier to use, thereby making it easier for young and old to brush their teeth. They are then subtly encouraged to brush their teeth more often. As a result automatic toothpaste dispensers promote oral hygiene.\nWhile automatic toothpaste dispensers have made it easier to dispense toothpaste onto a brush they are still not user-friendly for people who have arthritis in their hands or who have other physical limitations such as Parkinson's disease. These people have difficulty in holding the bristle end of a toothbrush under the dispensing spout of the automatic toothpaste dispensers and the toothpaste too often does not go onto the bristles but, rather, onto the counter or floor.\nWhether toothpaste tubes, manually operated toothpaste dispensers or automatic toothpaste dispenser, they all have a common problem. A person using them often touches the opening where the tooth is dispensed from such devices with the bristles of their toothbrush. This creates a hygienic problem because the bristles and handles of virtually all toothbrushes have germs, bacteria and sometimes viruses thereon. A fast rinse of a toothbrush under a running faucet after brushing does not properly clean a toothbrush. They still have germs, bacteria and viruses thereon. Very few families will have a separate toothpaste dispensers for each person in a family and germs, bacteria and viruses are thereby easily spread between members of a family or other group who share one of the many types of toothpaste dispensers described above. This problem is exacerbated when people with physical limitations, such as Parkinson's disease, must hold their toothbrush bristles against the opening where the toothpaste is dispensed in order to get the toothpaste on their toothbrush.\nTherefore, there is a need for means to dispense toothpaste using a dispenser that is shared by people while minimizing the transfer of germs, bacteria and viruses between the people. In addition, there is a need in the art for means to help a person position their toothbrush to receive toothpaste only onto their toothbrush and thereby minimizing any messes caused during dispensing toothpaste."}
{"text": "The peroxisome proliferator-activated receptors act as nuclear transcription factors and occur in three different forms, PPAR-alpha, PPAR-delta and PPAR-gamma. When activated, these receptors regulate the expression of genes involved in a variety of metabolic pathways. PPARα-regulated target genes are centrally involved in the metabolism of fatty acids. Specific PPARα effects include an increase in fatty acid oxidation in peroxisomes and mitochondria. Additional evidence indicates that PPARα activation leads to a variety of anti-inflammatory effects. Known activators include fibrates, in use clinically, as well as herbicides and phthalate plasticizers (Reddy, et al., Crit. Rev. Toxicol. 12:1-58 (1983)).\nThe PPAR-alpha receptor has been isolated and is known to be activated by fatty acids (Isseman, et al., Nature 347:645-650 (1990); Gottlicher, et al., Proc. Nat'l Acid. Sci. USA 89:4653-4657 (1992); Tugwood, et al., EMBO J. 11:433-439 (1992); Bardot, et al., Biochem. Biophys. Res. Comm. 192:37-45 (1993); Muerhoff, et al., J. Biol. Chem. 267:19051-19053 (1992); Marcus, et al., Proc. Nat'l Acid. Sci. USA 90:5723-5727 (1993)). There have been suggestions that activators of this receptor may be useful in treating skin conditions (U.S. Pat. No. 6,060,515) and obesity (U.S. Pat. No. 6,028,109). In addition, at least one report has suggested that activators of the gamma form of PPAR may be useful in treating inflammatory skin conditions (U.S. Pat. No. 5,981,586). A better understanding of the way in which the PPAR receptors are activated, especially in terms of endogenous pathways, may lead to new and better approaches to therapy."}
{"text": "Recycling of waste materials is highly desirable from many viewpoints, not the least of which are financial and ecological. Properly sorted recyclable materials can often be sold for significant revenue. Many of the more valuable recyclable materials do not biodegrade within a short period, and so their recycling significantly reduces the strain on local landfills and ultimately the environment.\nTypically, waste streams are composed of a variety of types of waste materials. One such waste stream is generated from the recovery and recycling of automobiles or other large machinery and appliances. For examples, at the end of its useful life, an automobile is shredded. This shredded material is processed to recover ferrous and non-ferrous metals. The remaining materials, referred to as automobile shredder residue (ASR), which may still include ferrous and non-ferrous metals, including copper wire and other recyclable materials, is typically disposed of in a landfill. Recently, efforts have been made to further recover materials, such as non-ferrous metals including copper from copper wiring and plastics. Similar efforts have been made to recover materials from whitegood shredder residue (WSR), which are the waste materials left over after recovering ferrous metals from shredded machinery or large appliances. Other waste streams that have recoverable materials may include electronic components (also known as “e-waste” or “waste electrical and electronic equipment (WEEE)), building components, retrieved landfill material, or other industrial waste streams. These recoverable materials are generally of value only when they have been separated into like-type materials. However, in many instances, no cost-effective methods are available to effectively sort waste materials that contain diverse materials. This deficiency has been particularly true for non-ferrous materials, and particularly for non-metallic materials, such as non-ferrous metals, including copper wiring. For example, one approach to recycling wiring has been to station a number of laborers along a sorting line, each of whom manually sorts through shredded waste and manually selects the desired recyclables from the sorting line. This approach is not sustainable in most economics since the labor component is too high.\nWhile some aspects of ferrous and non-ferrous recycling has been automated for some time, mainly through the use of magnets, eddy current separators, induction sensors and density separators, these techniques are ineffective for sorting some non-ferrous metals, such as copper wire. Again, labor-intensive manual processing has been employed to recover wiring and other non-ferrous metal materials. Because of the cost of labor, many of these manual processes are conducted in other countries and transporting the materials adds to the cost.\nMany processes for identifying and separating materials are know in the art. However, not all processes are efficient for recovering non-ferrous metals and the sequencing of these processes is one factor in developing a cost-effective recovery process. Also, many processes are “wet,” that is, they involve using water or other liquid media. Wet processes tend to be less cost effective, in part, because of the extra processing required to manage and dry materials.\nIn view of the foregoing, a need exists for cost-effective, efficient methods and systems for recovering materials from a waste stream, such as materials seen in a recycling process, including non-ferrous metals, in a manner that facilitates revenue recovery while also reducing landfill and, preferably using a dry process."}
{"text": "As a semiconductor memory device for storing a data which is excellent in a portability, there has been widely used a flash memory corresponding to a nonvolatile memory. A cost per bit of the flash memory has rapidly come down year by year, and a falling degree is sharper than a falling degree which is expected only by reduction of memory cell size. This has been achieved by reducing an area per bit on the basis of an artifice of an element structure or an introduction of a multilevel storage. In the semiconductor memory device, an area of a memory cell is reflected on a chip area, thereby being reflected on a manufacturing cost of the semiconductor memory device. A rough standard of the cost can be grasped by calculating how many times F2 a unit memory cell can be constructed, in which reference symbol F denotes a patterning dimension. At the present day, NAND type flash memory and AG-AND type flash memory can achieve a cell size which is approximately equal to a memory cell size corresponding to an ideally small 4F2. To make a memory cell smaller, it is necessary to line up the structure at a pitch smaller than 2F corresponding to twice of the lithography dimension F. To realize this, there can be considered such an artifice as to form a half pitch structure by utilizing a side wall of the once formed structure. As an example of the technique mentioned above, there can be listed up Japanese Patent Application Laid-Open Publication No. 08-055908 (Patent Document 1). And, there can be considered a technique of relaxing a pitch of a local data line (a local bit line) so as to connect to a global data line (a global bit line), however, in the case of relaxing the pitch of the local data line so as to connect to the global data line (the global bit line), it is necessary to set a select transistor to multi stages. As an example of a method of achieving a multi-stage select transistor with a small area penalty, there can be listed up Japanese Patent Application Laid-Open Publication No. 2004-31448 (Patent Document 2)."}
{"text": "Hydroxyl terminated polyethers find uses as raw materials of polyurethane and polyester resins. Others uses include cosmetics, plasticizers, surfactants and raw materials of these products. Conventional polyethers are terminated with primary or secondary alcoholic groups, the presence of which normally contributes to their desirablity for intended uses. It has been experienced in certain cases that the primary or secondary alcohol function is too active and therefore the presence of which is not desirable. For example, polyethers having too active an alcoholic function can lead to a premature reaction when reacted with polyisocyanates. Attempts have been made to remove the hydroxyl function by blocking with a group such as methyl or acyl, or replacing the hydroxyl group with chlorine. Obviously, polyethers having blocked or chlorinated terminals are not suitable for uses in which the presence of a hydroxyl function is essential, e.g. for the production of polyurethane and polyester resins."}
{"text": "For circuits such as switching power supply circuits and inverter circuits, semiconductor devices such as switching devices and diodes are used. Such semiconductor devices are required to have high breakdown voltage and low on-resistance. There is a trade-off relationship determined by device materials between the breakdown voltage and the on-resistance.\nWith the advancement in development of technology, low on-resistance for semiconductor devices has been realized down to the limit of silicon which is a major device material. In order to further improve the breakdown voltage and further reduce the on-resistance, it is necessary to change the device materials. By using GaN-based semiconductors such as GaN and AlGaN as switching device materials, it is possible to improve the trade-off relationship determined by the materials and to achieve rapidly-progressive high breakdown voltage and low on-resistance.\nHowever, for example, in a switching device using a GaN-based semiconductor, there is a problem of “current collapse” in which on-resistance increases when a high drain voltage is applied."}
{"text": "Systems for rapid thermal processing of glass panels for flat panel displays are known in the art. Such systems include a heating assembly and a transport assembly within an enclosure. Glass panels having one or more coatings enter the system in a horizontal orientation and are transported through the system on rotating quartz cylinders. The glass panels pass through a series of preheating zones at increasing temperatures. The preheating zones contain infrared heaters that are oriented perpendicular to the direction of movement of the glass panels. The panels then move through a region where they are exposed to intense ultraviolet radiation supplied by a xenon arc lamp mounted in a reflector. The panels then move through a cooling zone and exit the system at the end opposite the first preheating zone.\nAn important parameter in the operation of rapid thermal processing systems is temperature uniformity across the surface of the glass panel or other substrate. Prior art systems have provided generally satisfactory performance but have had certain drawbacks.\nThe system requires a non-contact temperature measurement technique. An optical pyrometer, which measures temperature based on sensed infrared energy, is typically utilized. The conventional prior art approach is to mount a pyrometer sensing head in a fixed position such that a line profile at one location along the direction of workpiece motion is produced. This configuration has limited temperature measurement capability and thus limits the ability to achieve uniform heating.\nPrior art rapid thermal processing systems typically utilize long filament infrared heaters positioned at 90.degree. to the direction of glass panel movement. This arrangement produces excellent longitudinal uniformity, but limited lateral uniformity due to edge effects. It is therefore desirable to provide a lamp configuration having improved lateral heating uniformity. Lateral uniformity refers to uniformity in a direction perpendicular to the direction of substrate movement.\nAs noted above, glass panels typically enter the rapid thermal processing system at one end and exit at the opposite end. The distance between the entrance end and the exit end increases as the dimensions of the glass panels being processed increases. This distance may be unacceptable in some fabrication facilities.\nA method for selectively heating a film on a substrate, such as a glass substrate for a flat panel display, is disclosed in U.S. Pat. No. 5,073,698 issued Dec. 17, 1991 to Stultz. A semiconductor wafer heating chamber, including an optical element between a light source and a wafer and an optical pyrometer for measuring wafer heating, is disclosed in U.S. Pat. No. 4,755,654 issued Jul. 5, 1988 to Crowley et al. A long arc lamp for semiconductor wafer heating is disclosed in U.S. Pat. No. 4,820,906 issued Apr. 11, 1989 to Stultz.\nAll of the known prior art rapid thermal processing systems have had one or more of the drawbacks discussed above. Accordingly, there is a need for improved methods and apparatus for rapid thermal processing of large area substrates."}
{"text": "Retail shopping facilities often have thousands of different items that can be purchased. In some instances, one or more of these items are uniquely identified through a bar code, radio frequency identifier (RFID) tag, or other identifier. When RFID tags are used, the retail facility uses RFID tag readers to detect and read the RFID tags. Often these RFID tag readers are positioned a distance from the tags when read. There is a need for improvements in RFID tag usage.\nElements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions and/or relative positioning of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present invention. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments of the present invention. Certain actions and/or steps may be described or depicted in a particular order of occurrence while those skilled in the art will understand that such specificity with respect to sequence is not actually required. The terms and expressions used herein have the ordinary technical meaning as is accorded to such terms and expressions by persons skilled in the technical field as set forth above except where different specific meanings have otherwise been set forth herein."}
{"text": "Computing devices, such as personal computers, mobile smartphones, personal digital assistants, tablets, televisions, digital glasses, video game systems, digital projection systems, and holographic display systems, can be used to view various kinds of computer stored data such as computer generated documents stored on digital media, scanned documents generated on external sources and then digitized and stored on digital media, and photographic images captured on various devices and then stored on digital media. Displaying various kinds of documents on computing device screens or on printed media may involve disparate and incompatible operating systems (OSs), which would require native applications specific to each OS, and sometimes specific to each device, in order to interact with the digitally stored data in a visual format.\nThis OS incompatibility can result in end users having to buy separate applications for each different category of devices, e.g., an application specific for a PC plus an application specific for a mobile device. In this regard, to access and manipulate a document across different device platforms, an end user must buy or download an application for each category of device utilized. In some cases, a user is required to buy or download two applications for the same category of device, e.g., an application specific for an Android device and an application specific for an iOS smartphone. While some envision mobile devices replacing desktop computing at some point in the future, the average consumer today continues to make use of various computing devices with different operating systems, such as personal computers (laptops and desktops), video game systems (Sony's PlayStation or Microsoft's Xbox) and other proprietary devices with native OSs. This situation also remains a burden to software developers, who, in order to address the widest potential market for their product, must spend the time, resources and money to develop their same basic application to run separately and specifically on each different device or platform.\nWith the advent of cloud computing, an extra wrinkle to the already challenging equation of accessing digitally stored data over the Internet has been created—that of efficiently retrieving and storing digital data on the cloud.\nApple's iOS and Google's Android OS currently dominate the world of mobile device OS, followed by Nokia, RIM's Blackberry, Microsoft's Windows Mobile and others. Although each mobile OS is different from an operations standpoint, they all share a common trait—their proprietary nature from both an internal architecture, as well as, the external applications (Apps) required for consumers to interact with the devices running on the OSs.\nA huge drawback of the current state of mobile computing is the lack of computational support from the server side. Although mobile devices have increased their capacity to serve and process data in the last couple of years, computing large tasks on mobile devices is still very costly; and it often requires cellular/wireless carrier subsidies to make the products affordable on the front-end and then through ongoing enhancements, new feature sets, improved processing power and other design feature or operational improvements, it creates frequent upgrade cycles."}
{"text": "A heat pump is understood to be a system involving a heat-circulation process of the type wherein a liquid cooling means at +2.degree. C. is evaporated in an evaporator and consequently absorbs the heat from a medium, for example water or air, at a temperature of +10.degree. C. The cooling-medium vapor of +2.degree. C. absorbs here approximately 10 kW. heating efficiency and is then drawn away by a compressor. The cooling-medium vapor is condensed to 15.5 bar in said compressor, is heated to +60.degree. C. and is fed to a condenser while absorbing a further 5 kW. (compressor output). Here the cooling-medium vapor emits again the absorbed 15 kW. at an unchanged high pressure and becomes liquid. The heating water, for example for a heating system of a house of the conventional type, is guided through the condensor and absorbs the heat which is emitted from the cooling-medium vapor, which heat is fed now in the heating water to the actual heating purpose.\nThe liquid cooling means at +60.degree. C. moves from the condensor on to an expansion valve, expands here and assumes a lower pressure of 3.5 bar at a temperature of +2.degree. C. This cooling means is returned to the evaporator and a new cycle of said means in the heat pump starts. (The aforementioned values are only exemplary values and are provided to aid one to easily understand the circulation processes of the heat pump. The physical principles of a heat pump are discussed in detail for example in \"VDI-Statusbericht Warmepumpe\" VDI-Verlag, GmbH, Dusseldorf 1976.).\nThe purpose of the present invention is now to provide a drive unit which is coupled with the compressor in such heat pumps, which drive unit can be operated and serviced easily and without any errors by nontrained persons, for example housewives. This purpose is attained according to the invention by the unit including a fuel-injection internal combustion engine which drives the compressor and a starter associated with said engine being fed power from an electric power line, wherein the starter can be connected to the electric power line by means of a switch.\nIn the preferred embodiment of the invention, the starter is constructed as an electromagnetic drive motor having a rotor and a stator. The rotor is secured to the crackshaft of the engine and serves at least partially as a flywheel. The stator has coil windings of which current from a three-phase a.c. power-line is fed.\nThe starting switch can be operated manually in heat pumps of a simple construction. In systems which have a high ease of operation, however, automation will be used, for example a thermostatically controlled relay, which, when falling below a fixed minimum value in the temperature of the heating water for the heating system of the house, closes the starting switch. This automation is preferable in view of the fact that a heat pump is used often only as an auxiliary heating system in a house, which is heated with heating oil, and must be started and turned off many times during changing weather temperatures.\nA preferable further development of the invention results when a coupling is arranged between the crankshaft of the engine and the driven shaft of the compressor, which coupling is constructed as a centifugal clutch. The compressor is connected in this manner to the driving fuel-injection internal combustion engine only after same has reached the necessary speed and the predetermined torque.\nFor the purpose of designing the drive unit as an easily transportable and installable unit, it is preferable to provide a common support frame, on which, with the interpositioning of vibration-damping elements therebetween, the structural components (internal combustion engine, starter, compressor) are secured. An arrangement which is very advantageous with respect to space results when the starter is arranged between internal combustion engine and compressor such that its housing is secured on the one side to the internal combustion engine and on the other side to the compressor."}
{"text": "Virtual reality (VR) is a computer technology that uses VR devices, such as headsets, sometimes in combination with physical spaces or multi-projected environments, to generate realistic images, sounds, and other sensations that simulate a user's physical presence in a three-dimensional (3D) virtual environment and allow the VR user to interact with the virtual environment. Many applications such as gaming, content consuming, and productivity have been developed using VR technology to provide user a more immersive experience.\nIn some cases, when a user needs to select a virtual element (for example, a virtual key) in a VR user interface (UI) to interact with a VR scenario, the user can select the virtual element in the VR UI by using an external device connected to the VR device. The external device connected to the VR device can be a smartphone, a joystick, or other handheld input devices. When using a VR application, the user can use the external device to control an operational focus (for example, a cursor) in the VR UI, move the operational focus to a virtual element, and select the virtual element in the UI by clicking a button or touchscreen on the external device.\nIn some cases, a clicking action can be simulated using a “hover timeout” event to select the virtual element in the VR UI. That is, a virtual element can be selected by hovering the operational focus on the virtual element for a pre-determined duration (for example, a three second timeout). When the operational focus stays on the virtual element for the duration, the virtual element is selected. Because the user needs to wait for a period of time to select the element, this method may not be suitable for VR applications that require real-time responses."}
{"text": "This invention relates to fuel elements for nuclear reactors, and more particularly to braces, spacers, or support lattices which are usually placed at predetermined distance(s) along the length of the fuel elements in order to provide lateral bracing and spacing, and to maintain the fuel elements in fixed positions.\nIn a nuclear reactor, the reactor core contains nuclear fuel which is typically in the form of fuel rods grouped together in fuel assemblies. The fuel assemblies are usually mechanically identical and can therefore, be interchangeable. Groups of fuel assemblies are arranged into a matrix to form a core capable of controlled fission reaction.\nEach fuel rod is typically a long member approximately 0.4 inches in diameter and 8 to 15 feet long containing fuel usually in the form of a stack of fuel pellets which are surrounded by tubular cladding. The fuel rods which make up an assembly are grouped together to form a plurality of longitudinally extending members which are supported vertically by two parallel end plates, an upper and a lower tie plate. These plates are usually connected to one another by tie rods, or other structural elements.\nEach fuel assembly or bundle may also include non-fuel bearing members. Examples include guide tubes to form passageways for control rods which assist in controlling the rate of fission, instrumentation tubes for in-core instrumentation, spacer capture rods, and water rods to modify the neutron moderation in the assembly. The spaces between adjacent fuel rods create flow channels through which coolant and/or moderator can circulate. In light water reactors, the coolant and moderator is water. Lateral bracing and spacing of the fuel rods in the fuel assembly are provided by spacers or spacer grids.\nThe fuel assembly or bundle, whether in a pressurized water reactor, boiling water reactor, high temperature gas cooled reactor, or any other type of reactor, functions in part to maintain the fuel rods in a fixed position, ideally free of vibration and restrained from bowing or other movement during normal and other operating conditions. In addition, by maintaining the fuel rods in fixed positions, proper cooling and neutron moderation can be achieved. Devices that assist in maintaining the fuel rods in fixed positions in the fuel assembly or bundle and which thereby facilitate proper fuel cooling are spacers.\nSpacers or spacer grids which provide lateral bracing are typically designed to allow differential axial expansion of the fuel rods. Springs incorporated in the spacer grids are most frequently used to permit some sliding of the fuel rods with respect to the spacer grids. In some of the designs, the spacer grid is free to move axially a small amount to accommodate minor changes in the axial length of the fuel rods during irradiation.\nIf spacers wee to be rigidly connected to the fuel rods as well as to structural members of the fuel assembly, then relative axial movement due to rod growth and thermal expansion of adjacent rods can cause local fuel rod skewing and bowing.\nBy being positioned at regular intervals, spacers maintain rod-to-rod spacing along the length of the fuel assembly. Spacers are typically made of zirconium based alloy sheet material or sometimes from Inconel or stainless steel, and are built up from a relatively large number of different intricately shaped strips that are fitted together by hand and subsequently welded. Sometimes, short sections of tubing are used that are welded to one another along parts of their edges. The spacers have an egg crate shape and each spacer cell includes dimples and/or springs to maintain the desired rod-to-rod spacing. Thus, the springs and dimples keep the fuel rods in their proper lateral positions. But, under the influence of irradiation, the springs are prone to relax and this can lead to undesirable changes in fuel rod pitch (i.e. rod-to-rod spacing) or it may cause gaps or spaces to develop between fuel rods and the springs and dimples, and increases the likelihood that the rods and/or spacer grids will vibrate. Such gaps, changes in fuel rod pitch, and vibration may lead to fuel rod fretting and failure.\nSprings necessarily contact the fuel cladding and can cause fretting even if they are not subject to relaxation. Furthermore, as the fuel is irradiated, the fuel rods undergo a shrinkage or diameter reduction known as \"creepdown\" which can result in gaps between the fuel rod cladding and the springs or dimples which in turn can cause or contribute to fuel rod fretting.\nSince even the slightest repeated relative movement between contacted fuel rods and spacers can result in fretting of the fuel cladding, minimization of such relative movement is desired. Thus, loads due to thermal expansion of fuel rods, particularly fuel cladding, and associated tensile loads from such expansion must be considered in spacer design.\nSpacers should be thin members and have minimal cross-sectional area. Ideally, they are invisible to moderator and coolant flow. Spacer designs which reduce the flow area, also increase flow resistance and restrict coolant flow causing undesirable pressure drops. Thus, the particular physical configuration of a spacer can create or contribute to local or even non-local undesirable flow redistribution, restriction, or distortion.\nTypically, the fabrication of spacers requires extensive labor in shaping the separate parts and in assembling and welding these parts to form a spacer grid. Many of these operations can be automated. However, even with automated spacer fabrication, assembly and welding, extensive labor is required in inserting the fuel rods through the spacers to form a fuel bundle. Even with simplified spacer designs, the assembly of spacer and fuel rods is intricate, time consuming and costly. Whenever spacer designs include springs and/or dimples, even non-intricate assembly techniques become particularly labor intensive. Labor costs can be substantial, particularly in light of the rigorous quality control standards that are applied to nuclear reactor components. When spacer design and fuel rod-spacer assemblies include spring or dimple elements, assembly of the fuel rods into a fuel bundle often result in scratching or damaging the fuel rod cladding surface which may lead to rejection of fuel assemblies during pre-operational quality control tests and inspections and/or fuel rod cladding failures during reactor operation.\nIt would thus be an advantage if spacers did not include springs and dimples or similar contact devices which relax from irradiation.\nIt would be an additional advantage to limit the repeated relative movement between spacers and the contacted fuel rod cladding.\nIt would be a further advantage if the actual physical connection between each spacer and fuel rod was not mechanical but was welded, fused, or bonded thereby avoiding the possibility of fretting of the fuel cladding by its spacer.\nIt would be a further advantage if the effective coolant flow area surrounding each fuel rod is maximized by minimizing the spacer cross sectional surface area which impedes coolant flow.\nIt would be yet a further advantage if the flow characteristics of a completed fuel assembly with fuel rod spacers could be improved by providing effectively unhindered coolant (e.g. water and steam) flow along the length of the fuel rod between the fuel cladding and its spacer and unhindered coolant flow within the area in the center of adjacent spacers by keeping the center free from any restrictions.\nIt would be a further advantage to minimize radial changes in the spacer structure in order to minimize changes in the pitch of the fuel rods which may affect the neutronic and thermal performance of the assembly.\nIt would be yet another advantage if spacers could be made of materials having low neutron cross section as well as having small mass.\nIt would be an advantage if each fuel rod was connected to its own spacer which remained integral with the fuel rod.\nIt would be an advantage if each fuel rod/spacer was mechanically identical.\nIt would be an advantage if the spacers were not rigidly connected to one another but were left free to move in the axial direction and thereby accommodate rod growth and thermal expansion differences between adjacent fuel rods.\nIt would also be an advantage if the spacer is not connected to the fuel assembly tie or spacer capture rod or to the guide tubes or other structure members so that axial movement of the rods is unrestricted and axial forces are not exerted on the rods thereby minimizing fuel rod bow.\nIt would also be an advantage if a portion of the spacer could deform in order to accommodate the reduction in diameter of the fuel rods due to creep without changing the distance between the spacer and fuel rod cladding as well as the distance between adjacent spacers.\nIt would be yet a further advantage if the fuel rod surface would not be exposed to any spacer parts that may cause fretting and if any fretting is to take place, that it would occur on the spacers and not the fuel rods.\nIt would be an additional further advantage if a spacer having each of the above advantages could be fabricated, tested and inspected at a lower cost than conventional spacers and at the same time improving overall quality and reliability.\nIt would be another advantage to have spacers which could be attached to and form a sleeve around fuel rods and avoid the prior art method of inserting the fuel rods through and interlacing with the spacer grids.\nIt would be yet a further advantage if each fuel rod and its integral spacer was such that their assembly into a fuel bundle merely involved stacking of individual fuel rods one on top of another without scratching any of the fuel rods."}
{"text": "The present invention relates generally to a computer implemented method for selecting responsive actions. More specifically, the present invention relates to referencing earlier attempted actions, if sufficiently current, and responding with at least one alternate action.\nSystems management software installed on data processing systems are used to detect and catalog events. An event is a signal that indicates a condition of a resource of a data processing system. An event can be, for example, a condition related to a rotation rate of a fan. Accordingly, the event can be a signal generated to identify both the resource and the condition of the resource. Some events can be asynchronously generated from an error condition. In addition, multiple events may be aggregated together to form an event summary. An event summary is an event that describes at least two other events. For example, an event summary can be a count of the number of times a fan speed exceeded a threshold in the previous hour. A resource is a physical data processing hardware or virtual resource, backed by physical data processing hardware. A resource can correspond to a system, subsystem, software program, or device. A resource can have attributes such as a temperature. An event, on the other hand, may be a request such as a hypertext protocol request at a server, or a report that a processor on that server has met or exceeded a temperature. Accordingly, the use of events can be helpful to collect information about a data processing system's status.\nOne or more events can be indicative of an underlying problem or anomaly. Accordingly, it is helpful that system administrators pay attention to some events under some conditions."}
{"text": "The present invention relates to the flavor industry. It concerns more particularly a process for producing a washed citrus oil flavor, which comprises treating a citrus oil with a mixture of propylene glycol and an alkanol from C1 to C4. The invention also relates to clear beverages flavored with citrus washed oils as formulated by the process of the invention.\nThe importance in the flavor industry of various citrus oils has resulted in intensive research directed towards processing methods, compositional analyses, quality and utilization.\nDetailed compositional data concerning the commercially important citrus oils can be found in several reviews. Extensive lists of various chemical compounds identified as components of orange, grapefruit, tangerine, lemon, and lime oils have been described.\nThe flavoring characteristics of the citrus oils are a function of the properties of certain components identified in the oils. Furthermore, d-limonene, which is the main component present in citrus oils, acts as a carrier of flavor for other oxygenated compounds mainly responsible for the citrus flavor characteristics.\nThe extraction from a citrus oil of the oxygenated compounds which are responsible for the citrus flavor is therefore a field of huge interest for the flavor industry. It is well known in the art that flavoring agents may be prepared by extracting from folded oils the oxygenated compounds desired in an extracted or washed flavor profile.\nThe formulation of washed oil flavors disclosed in the prior art uses water or propylene glycol as solvents for the extraction of useful flavor compounds from a citrus oil. For instance, U.S. Pat. No. 3,809,757 describes the preparation of a flavoring extract from orange oil by selective extraction, using propylene glycol and a hydrocarbon solvent.\nThe quality of an extract or a wash is mainly characterized by the selectivity of the extracted compounds. The scope sought in the flavor industry when formulating a wash from a citrus oil, is therefore to selectively obtain the key citrus flavor compounds i.e., the oxygenated compounds such as neral, geranial, neral acetate, geranial acetate, decanal, nootkatone, etc. Other requirements such as clearness or flash point requirements are necessary for obtaining an optimized washed oil product.\nNow, we have been able to find a novel process for the preparation of a citrus washed oil flavor, combining propylene glycol usually utilized for the extraction of oxygenated compounds, with an alkanol from C1 to C4. This process makes it possible to obtain a surprisingly selective wash of the key flavor compounds present in a citrus oil, which further meets all the requirements of quality and security needed for such a product. The present invention is thus related to a process for producing a citrus washed oil flavor, which comprises treating a citrus oil with a mixture of propylene glycol and an alkanol from C1 to C4. In a preferred embodiment the alkanol is selected from the group consisting of ethanol and isopropanol.\nThe process of the invention advantageously allows to obtain citrus washed oil flavors which meet the clarity, flash point and usage level objectives required in the art. The washes obtained by the process according to the invention have a flash point greater than 100xc2x0 F. (38xc2x0 C.) (i.e. they belong to the Class I liquids), cut clear in water and have a light, clean profile with the hint of fruit flavor from 0.05 to 0.2%, preferably at a 0.1% use level. Moreover, analyses of the washed citrus oil flavors formulated respectively with and without the use of an alkanol from C1 to C4, clearly indicate that washes produced with the use of some alkanol are superior to those made with propylene glycol on its own. This result will be demonstrated in the comparative tables below, in the case of ethanol.\nUnexpectedly, and besides the fact that alkanols from C1 to C4 suffer from the standpoint of their low flash point, we have now been able to optimize the formulation of citrus washed oils using the alkanol at the right level to extract the oxygenated compounds while meeting the flash point limitations.\nThe proportion of alkanol used in the extracting mixture with propylene glycol can vary in a certain range of values. Typically it can be comprised between 3 and 10% and is preferably comprised between 5 and 7% of the total weight of the composition comprising the citrus oil, propylene glycol and the alkanol.\nAs mentioned above, among the alkanols from C1 to C4, ethanol and isopropanol are preferred.\nThe starting product for the process of the invention may be any commercially available citrus oil, in the form of a cold pressed oil, or a folded oil, for instance. Non limiting examples of citrus oils which can be used for the process of the invention include lemon, berry, orange, grapefruit, tangerine, lime, kumquat, mandarin, bergamot and other essential oils or a mixture thereof.\nThe selectivity of the extraction system may be demonstrated by the comparison between citrus washes made with some ethanol, and washed oil flavors made with 75% propylene glycol: 25% flavor oil, i.e. ethanol free.\nThe following tables compare the washes obtained with propylene glycol combined with ethanol, and those obtained with propylene glycol on its own.\nTable 1 reports the results for a lemon wash while Table 2 reports the results for a grapefruit wash.\nIt can be easily concluded from the results reported in the above tables that the use of some ethanol is advantageous in extracting the key flavor compounds (i.e. oxygenated compounds such as neral, geranial, neral acetate, geranial acetate, decanal, nootkatone, etc.) from a citrus oil.\nThe invention also relates to a method to impart, improve, enhance or modify the flavoring properties of a flavoring composition or a flavored product, which comprises adding to said composition or product a citrus washed oil flavor produced by the process of the invention.\nAnother object of the invention is a flavoring composition or a flavored product such as a clear beverage, comprising a citrus washed oil flavor as an active flavoring ingredient, the proportion of the latter being comprised between 0.05 and 0.2% of the total weight of the composition or product."}
{"text": "1. Field of the Invention\nThis invention relates generally to memory design. Particularly, this invention relates to a new cell for implementation in the Static Random-Access Memory (SRAM) thus taking advantage of the base current reversal phenomenon in a Complementary Metal Oxide Semiconductor (CMOS) for the construction of compatible high gain gated lateral Bipolar Junction Transistor (BJT).\n2. Description of the Prior Art\nThe conventional CMOS SRAM cell essentially consists of a pair of cross-coupled inverters as the storage flip-flop or latch, and a pair of pass transistors as the access devices for data transfer into and out of the cell. (A large number of different cell configurations are cited in the literature, e.g.,B. Prince, Semiconductor Memories--A Handbook of Design, Manufacture, and Application, 2nd ed., New York: John Wiley & Sons, Inc., 1991). Thus, a total of six Metal Oxide Semiconductor Field Effect Transistors (MOSFETs) or four MOSFETs plus two very high resistance load devices are required for implementing a conventional CMOS SRAM cell. However, to achieve high packing density, it is the usual practice to reduce the number of the devices needed for realizing a CMOS SRAM cell or the number of the devices for performing the Write/Read operation. Especially for the case of very high resistance load devices, increased process complexity, extra masks, and high fabrication cost are required for forming the undoped polysilicon layers or the Thin Film Transistor (TFT) on the oxide and thus saving the chip area; however, the corresponding product yield is not high. Therefore, more efforts are needed to further reduce the areas occupied by the chip while improving the production yield.\nThe structures of the conventional SRAM are shown in FIGS. 1 and 2. FIG. 1 shows a circuit schematic of a conventional SRAM cell configuration. The cell comprises a pair of cross-coupled inverters, p-MOSFET 11 and n-MOSFET 13, and p-MOSFET 16 and n-MOSFET 14, as the storage flip-flop or latch. In each inverter the gates of p-MOSFET and n-MOSFET are tied together and connected to the output of another inverter. The output of each inverter are the drains of p-MOSFET and n-MOSFET which are tied together. The conventional cell employs a pair of pass transistors, n-MOSFETs 12 and 15, as the access devices for data transfer into and out of the cell. Two column lines DATA and DATA and two row select lines W and W are depicted. This conventional cell requires 6 MOSFETs.\nFIG. 2 shows a circuit schematic of a conventional CMOS cell using undoped polysilicon layer or thin-film transistor for providing very high resistance loads 20. FIG. 2 is very similar to FIG. 1 except for the type of the load. This cell requires complex processes and the resulting yield is low. This conventional cell requires four MOSFETs formed on the silicon plus two undoped polysilicon layers or TFTs formed on the oxide.\nA Bipolar-CMOS (BICMOS) process has recently been introduced to achieve the above-mentioned object (refer to the paper: K. Sakui, et al. \"A new static memory cell based on reverse base current (RBC) effect of bipolar transistor,\" IEEE IEDM Tech. Dig., pp. 44-47, December 1988). In this BICMOS process, only two devices are needed for a SRAM cell: one vertical bipolar transistor and one MOSFET as a pass device. However, extra processing steps and increased masks are required along with special deep isolation techniques, resulting in high fabrication cost and process complexity. Also, the yield of the SRAM products utilizing such complex BICMOS process is usually low compared with the existing standard CMOS process. Therefore, the SRAM products fabricated in a complex BICMOS process cannot provide sufficient competition to the conventional CMOS SRAM products.\nVery recently, a new phenomenon of base current reversal has been observed in a CMOS compatible high gain gated lateral bipolar transistor (refer to the paper: Tzuen-Hsi Huang and Ming-Jer Chen, \"Base current reversal phenomenon in a CMOS compatible high gain n-p-n gated lateral bipolar transistor,\" IEEE Trans. Electron Devices, Vol. 42, No. 2, pp.321-327, February 1995). This new phenomenon has been extensively investigated and has been found to have many applications. This invention employs this new phenomenon in the existing standard CMOS process for realizing a new SRAM cell constructed by only two MOSFETs. This new SRAM also features only one-sided peripheral circuitry for Read/Write action. Therefore, this invention can efficiently save the chip area with high yield since it is fully compatible with the existing low-cost standard CMOS process."}
{"text": "The present invention relates to particles, such as carbon black, having an attached stable free radical which permits the formation of modified particles having blocked radical sources. The present invention further relates to methods of preparing and using such modified particles.\nThe role of carbon black as a thermal stabilizer in polymeric systems is discussed by W. L. Hawkins, R. H. Hansen, W. Matreyek, F. H. Winslow; J. Applied Polymer Science, vol. 1, pages 37-42, 1959; J. T. Cruver, K. W. Rollmann: J. Applied Polymer Science, vol. 8, pages 1169-83, 1964, and by G. Ivan, M. Giurgina: Revue Rotimaine de Chemie, vol. 29, number 8, pages 639-646, 1984. In each discussion the mechanism is through phenolic and quinone oxygen groups on the carbon black surface behaving either as radical traps or peroxide decomposers. One skilled in the art, however, would consider these sites unavailable as initiating sites for polymerization processes.\nBelmont et al. (J. A. Belmont, J. M. Funt: International Rubber Conference, Essen, Germany, June 24-27 1991) identified the presence of peroxide groups, typically in the range of 0.1 to 0.4 micromoles/sq meter, on the carbon black surface. However, the majority (greater than 80%) of the peroxide groups are thermally stable to heat treatment at 200xc2x0 C. and hence cannot be considered to be potential initiating sites for radical polymerization.\nTsubokawa et al. (K. Fujiki, N. Tsubokawa, Y. Sone: Polymer J., vol. 22, number 8, pages 661-70, 1990, and N. Tsubokawa: Prog. Polymer Science, vol. 17, pages 417-70, 1992) discuss growing polymers from an activated carbon black surface by first attaching a reactive group via the oxygen groups on the carbon black surface. Typical examples include the use of glycidyl methacrylate where the glycidyl group reacts with phenolic hydroxyl groups on the carbon black surface providing a vinyl functionality; the reaction of 4,4xe2x80x2 azo bis-(4-cyanovaleric acid) whereby the isocyanate group reacts with phenolic hydroxyl groups and subsequent heating decomposes the azo group to generate an alkyl radical; and the reaction of the phenolic hydroxyl groups with butyl lithium which can then be used as an initiation site for anionic polymerization.\nAll of these approaches require the presence of oxygen groups on the carbon black surface. Since the level of reactive hydroxyl and carboxylic acid groups, relative to the above approaches, present on a typical furnace or thermal carbon black is typically 0.01 to 0.1 micromoles/sq meter, the number of potential initiation sites is quite low.\nFurther, subsequent polymerization from these activated sites will most likely occur in the normal manner for free radical polymerization with the chains being irreversible terminated by chain combination reactions, combination with unreacted oxygen groups on the carbon black surface, and/or the addition of chain stoppers. In all cases the polymerization cannot be reinitiated. Accordingly, there is a need to provide particles with attached stable free radicals which overcome the above-described limitations.\nIn accordance with the purposes of the present invention, the present invention relates to a modified particle which includes a particle having an attached group having the formula: \nwherein A represents an aromatic group or an alkyl group; R1 represents a bond, an arylene group, an alkylene group, \nwherein R4 is an alkyl or alkylene group or an aryl or arylene group; R2 and R3, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR5, xe2x80x94NHR5, xe2x80x94NR5R5 or xe2x80x94SR5, wherein R5, which can be the same or different, represents hydrogen, an alkyl group or an aryl group; and SFR represents a stable free radical.\nThe present invention further relates to a modified particle or aggregate wherein the particle or aggregate is a carbon-metal multiphase aggregate, a carbon-silicon containing species multiphase aggregate, a metal oxide, or a metal hydroxide. Attached to the particle or aggregate is a group having the formula: \nwherein CoupA represents a Si-containing group, a Ti-containing group, a Cr-containing group, or a Zr-containing group; R8 and R9, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR10, xe2x80x94NHR10, xe2x80x94NR10R10, or xe2x80x94SR10, wherein R10, which can be the same or different, represents hydrogen, an alkyl group or an aryl group; SFR represents a stable free radical, and n is an integer of from 1 to 3.\nThe present invention further relates to a modified particle with an attached polymer or blocked radical sources, wherein the modified particle has an attached group having the formula: \nwherein X represents blocked radical sources or a polymer formed from at least one polymerizable vinyl or diene based monomer, A represents an aromatic group or alkyl group; R1 represents a bond, an arylene group, an alkylene group, \nwherein R4 is an alkyl or alkylene group or an aryl or arylene group; R2 and R3, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR5, xe2x80x94NHR5, xe2x80x94NR5R5 or xe2x80x94SR5, wherein R5, which can be the same or different, represents an alkyl group or an aryl group; and SFR represents a stable free radical.\nThe present invention also relates to another modified particle having an attached polymer. The particle is a carbon-metal multiphase aggregate, a carbon-silicon containing species multiphase aggregate, a metal oxide, or a metal hydroxide. Attached to the particle is a group having the formula: \nwherein CoupA represents a Si-containing group, a Ti-containing group, a Cr-containing group, or a Zr-containing group; R8 and R9, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR10, xe2x80x94NHR10, xe2x80x94NR10R10, or xe2x80x94SR10, wherein R10, which can be the same or different, represents an alkyl group or an aryl group; SFR represents a stable free radical, X is blocked radical sources or a polymer formed from at least one polymerizable vinyl or diene containing monomer, and n is an integer of from 1 to 3.\nThe present invention, in addition, relates to a method for preparing and using these various modified particles. The method includes reacting a particle having an attached vinyl substituted alkyl group with a reactive free radical source and a stable free radical source to form a reaction product. The present application further relates to forming the modified particles with an attached blocked radical source or polymer by reacting the reaction product with a polymerizable vinyl or diene containing monolayer.\nThe present invention further relates to a method of making a modified particle comprising a particle having attached a group having the formula: \nwherein A represents an aromatic group or an alkyl group; R1 represents a bond, an arylene group, an alkylene group, \nwherein R4 is an alkyl or alkylene group or an aryl or arylene group; R2 and R3, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR5, xe2x80x94NHR5, xe2x80x94NR5R5, or xe2x80x94SR5, wherein R5, which is the same or different, represents an alkyl group or an aryl group; and SFR represents a stable free radical, wherein said method comprises reacting a) a particle having attached an aromatic group or an alkyl group with a group containing an abstractable proton and which leaves behind a carbon-centered radical when the proton is abstracted with b) a reactive free radical source and c) a stable free radical source.\nThe present invention, in addition, relates to a modified particle which includes a particle having an organic group having a xe2x80x94SFR group, where the organic group is directly attached to the particle. The modified particle can also be a particle having an organic group having an xe2x80x94Xxe2x80x94SFR group, where the organic group is directly attached to the particle.\nThe present invention further relates to polymer-modified particles or aggregates where the xe2x80x94SFR group is replaced or terminated with a terminated moiety, for example, a hydrogen atom, a hydroxy group, or a halide group.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory only and are intended to provide further explanation of the present invention, as claimed.\nIn one embodiment, the present invention relates to a modified particle comprising a particle having attached a group having the formula: \nwherein A represents an aromatic group or an alkyl group which is attached to the particle; R1 represents a bond, an arylene group, an alkylene group, \nwherein R4 is an alkyl or alkylene group or an aryl or arylene group; R2 and R3, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR5, xe2x80x94NHR5, xe2x80x94NR5R5, or xe2x80x94SR5, wherein R5, which can be the same or different, represents an alkyl group or an aryl group; and SFR represents a stable free radical.\nThe particle, to be modified, can be any particle capable of having a group represented by any one of formulas described herein attached thereto, such as formula (I)-(IV). Preferably, the particle has a carbon reactive site (Cxe2x80x94), since in a preferred embodiment for the groups of formula (I) and (III), xe2x80x94SFR and xe2x80x94Xxe2x80x94SFR are preferably attached through a carbon bond on the particle. The particle, for instance, can be carbon products, colored pigments other than carbon products, metal oxides (e.g., silica), metal hydroxides, multiphase aggregates comprising a carbon phase and a silicon-containing species phase, multiphase aggregates comprising a carbon phase and a metal-containing species phase, and the like. The carbon may be of the crystalline and/or amorphous type. Examples of carbon products include, but are not limited to, graphite, carbon black, vitreous carbon, activated charcoal, activated carbon, carbon fiber, and mixtures thereof. Finely divided forms of the above are preferred. Most preferably, the particle is a carbon product, and most preferably carbon black.\nThe multiphase aggregate containing the carbon phase and the silicon-containing species phase can also be considered a silicon-treated carbon black aggregate and the multiphase aggregate containing a carbon phase and a metal-containing species phase can be considered to be a metal-treated carbon black aggregate as long as one realizes that in either case, the silicon-containing species and/or metal-containing species are a phase of the aggregate just like the carbon phase. The multiphase aggregates do not represent a mixture of discrete carbon black aggregates and discrete silica or metal aggregates and is not a silica coated carbon black. Rather, the multiphase aggregates used in the present invention include at least one silicon-containing or metal-containing region concentrated at or near the surface of the aggregate (but put of the aggregate) and/or within the aggregate.\nThe aggregate, thus contains at least two phases, one of which is carbon and the other of which is a silicon-containing species, a metal-containing species, or both. The silicon-containing species that can be a part of the aggregate is not attached to a carbon black aggregate like a silica coupling agent, but actually is part of the same aggregate as the carbon phase. For instance, when the multiphase aggregate having a carbon phase and a silicon-containing species phase is examined under STEM-EDX, the silicon signal corresponding to the silicon-containing species is found to be present in individual carbon black aggregates. By comparison, for example, in a physical mixture of silica and carbon black, STEM-EDX examination reveals distinctly separate silica and carbon black aggregates.\nThe metal-treated carbon blacks are aggregates containing at least one carbon phase and at least one metal-containing species phase. The metal-containing species include compounds containing aluminum, zinc, magnesium, calcium, titanium, vanditimi, cobalt, nickel, zirconium, tin, antimony, chromium, neodyitium, lead, tellurium, barium, cesium, iron, and molybdenum. Preferably, the metal-containing species phase is an aluminum-or zinc-containing species phase. The metal-containing species phase(s) can be distributed through at least a portion of the aggregate and is an intrinsic part of the aggregate.\nFurther, it is within the bounds of the present invention to have a metal-treated carbon black containing more than one type of metal-containing species phase or the metal-treated carbon black can also contain a silicon-containing species phase and/or a boron-containing species phase. For example, the metal-treated carbon black of the present invention can have an aggregate comprising a carbon phase, an aluminum-containing species phase, and a zinc-containing species phase. Accordingly, the multiphase aggregates used in the present invention can have two or more different types of metal-containing species phases and/or additional non-metal species phases.\nPreferably, for purposes of the present invention, the amount of elemental silicon and/or elemental metal present in the multiphase aggregate is from about 0.1 to about 25 wt. %, more preferably, from about 0.5 wt. % to about 10 wt. %, and most preferably from about 0.2 wt. % to about 5.4 wt. %, by weight of the aggregate.\nThe details of making the multiphase aggregates are explained in U.S. patent application User. Nods.: 08/446,141, filed May 22, 1995; 08/446,142, filed May 22, 1995; 08/528,895, filed Sep. 15, 1995; and 08/750,017, filed Nov. 22, 1996, which is a National Phase Application of PCT No. WO 96/37547, filed May 21, 1996, U.S. patent application Ser. No. 08/828,785, filed Mar. 27, 1997, U.S. patent application Ser. Nos. 08/837,493 filed Apr. 18, 1997 and 09/061,871 filed Apr. 7, 1998. All of these patent applications are hereby incorporated in their entireties herein by reference.\nA silica-coated carbon product can also be used as the particle, such as that described in PCT Application No. WO 96/37547, published Nov. 28, 1996, which is hereby incorporated in its entirety herein by reference.\nThe colored pigment is any pigment which can be modified with the attachment of an aromatic group or an alkyl group. The colored pigments other than carbon products include, but are not limited to, black, blue, brown, cyan, green, violet, magenta, red, yellow, as well as mixtures thereof. Suitable classes of colored pigments include, but are not limited to, anthraquinones, phthalocyanine blues, phthalocyanine greens, diazos, monoazos, pyranthrones, perylenes, heterocyclic yellows, quinacridones, and (thio)indigoids. Specific examples and further information concerning the colored pigments and methods of making colored pigments with attached aromatic groups or alkyl groups are described in PCT Publication No. WO 97/47699, incorporated in its entirety by reference herein.\nParticles useful for the present invention may, for example, have primary particles sizes in the general range of from about 10 nm to about 500 nm, and preferably from about 10 nm to about 250 nm, and primary aggregate sizes in the general range of from about 50 nm to about 100 microns, preferably from about 50 nm to about 10 microns, still more preferably from about 75 nm to about 1 micron. The BET surface area of these particles can be any suitable surface area and preferably ranges from about 10 m2/g to about 2000 m2/g and more preferably, from about 10 m2/g to about 1,000 m2/g, and still more preferably from about 50 m2/g to about 500 m2/g; and the particle structure preferably ranges from about 10 cc/100 g to about 1000 cc/g, more preferably, from about 50 cc/100 g to about 200 cc/100 g.\nThe number of SFR groups directly attached to the particle prior to polymerization can be any amount which can be effectively attached. For instance, the number of xe2x80x94SFR groups may range from about 0.01 mmole/g (of particle) to about 10 mmole/g, or from about 0.1 mmole/g to about 4 mmole/g, or from about 0.05 mmole/g to 4 mmole/g or from about 0.5 mmole/g to about 3 mmole/g, or from about 0.1 mmol/g to about 2 mmol/g.\nWhen the particle is a multiphase aggregate, like a particle comprising a carbon phase and a silicon-containing species phase, preferably the group of formula (I) or (III) is attached at least, if not exclusively, on the carbon phase.\nWith regard to the A group, which represents at least one aromatic group or an alkyl group, any aromatic group or alkyl group may be used. Unlike the polymerizable monomer discussed later, the aromatic group or the alkyl group is not a polymer and is not polymerized. Examples include, but are not limited to, arylene groups. Preferred arylene groups include, but are not limited to, phenylene and naphthalene groups. The aromatic group includes, but is not limited to, unsaturated cyclic hydrocarbons containing one or more rings. The aromatic group may be substituted or unsubstituted. Aromatic groups include aryl groups (for example, phenyl, naphthyl, anthracenyl, and the like), and heteroaryl groups (for example, imidazolyl, pyrazolyl, pyridinyl, thienyl, thiazolyl, furyl, triazinyl, indolyl, and the like). The alkyl group is preferably a C1-C12 alkyl group, and may be branched or unbranched, and is preferably ethyl.\nWith respect to R1, preferred arylene groups include, but are not limited to, benzene ring containing groups. Preferred alkylene groups include, but are not limited to, C1-C18 alkylene-containing groups. These groups can be linear, branched, or unsaturated. These examples of arylene and alkylene groups can also be considered examples of R4. Preferred alkyl groups for R4 are C1-C20 alkyl groups, more preferably C1-C5 alkyl groups and preferred aryl groups are phenyl, biphenyl, and naphthyl.\nWith respect to R2 and R3, which can be the same or different, examples of alkyl groups (e.g. C1-C20 alkyl group) include, but are not limited to, methyl, ethyl, propyl, butyl, and the like. Preferably, the alkyl group is a C1-C5 alkyl group. Examples of aryl groups include, but are not limited to, phenyl, biphenyl, and naphthyl. The alkyl and aryl groups mentioned here as well as the arylene and alkylene groups mentioned throughout can be unsubstituted or substituted for purposes of the present invention. R5 can be the same type of alkyl and aryl groups mentioned above with respect to R2 and R3.\nSFR, which is the stable free radical, can be any radical capable of capping the remaining portion of the group attached onto the particle. Examples of the SFR include, but are not limited to, nitroxide free radicals such as 2,2,5,5-tetramethyl-pyrrolidinlyloxy and 2,2,6,6-tetramethyl-piperindinyloxy, organic hydrazyl compounds, organic verdazyl compounds, organic aroxyl compounds (e.g., 2,4,6 tri-tertiary butyl phenoxy radical, galvinoxyl (2,6 ditertiary butyl alpha 3,5 ditertiary butyl oxo 2,5 cyclohexadiene-1 ylidene para tolyoxy) radical), aryl alkyl or aryl cycloalkyl where the unpaired electron is on a carbon atom, substituted triphelnyl methyl, substituted triphenyl amine, and derivatives of these compounds.\nThe SFR used in the present invention preferably has the formulas: \nor\nxe2x80x83xe2x80x94Oxe2x80x94Ar2xe2x80x83xe2x80x83(VI)\nwherein R6 and R7, which can be the same or different, each represent a substituted or unsubstituted alkyl or cycloalkyl group; or R6 and R7 can together form a cycloalkyl group; Ar2 represents a substituted or unsubstituted aromatic group. Representative examples of the alkyl and aromatic groups can be the same as described above with respect to the substitutents R2 and R3.\nAnother embodiment of the present invention relates to a modified particle or aggregate having a group of formula (II) attached thereto. The particle can be a metal oxide, a metal hydroxide, an aggregate comprising at least one carbon phase and at least one metal-containing species phase, or an aggregate comprising at least one carbon phase and at least one silicon-containing species phase. Attached to this particle or aggregate is a group having the formula: \nwherein CoupA represents a Si-containing group, a Ti-containing group, a Cr-containing group; or a Zr-containing group; R8 and R9, which can be the same or different, represent hydrogen, an alkyl group, an aryl group, xe2x80x94OR10, xe2x80x94NHR10, xe2x80x94NR10R10, or xe2x80x94SR10, wherein each R10, which can be the same or different, represents an alkyl group or an aryl group; SFR represents a stable free radical, and where n is an integer of from 1 to 3. Preferably, CoupA is attached to the particle or aggregate, especially in the case of a Si-containing group, via an oxygen radical which can form a part of a CoupA.\nExamples of Si-containing groups include, but are not limited to, dimethylsilylmethyl, dialkoxysilylmethyl, and the like. Examples of Ti-containing groups include, but are not limited to, alpha substituted tri-acetyl titanate and the like. Examples of Zr-containing groups include, but are not limited to, dialpha methoxy neopentylzirconate, aluminum zirconates, and the like.\nExamples of the substituents R8 and R9 can be the same as the substituents R2 and R3 mentioned above. Likewise, examples of the substituent R10 can be the same as the substituent R5 discussed above. Also, the SFR is the same as discussed above.\nThe modified particles having the attached group of the formulas, such as formula (I), can be made in the following manner. A particle, such as carbon black, can first have a vinyl substituted aromatic group or a vinyl substituted alkyl group attached to the particle. This attachment can be accomplished by the methods described in PCT International Application No. WO 96/18688 and U.S. Pat. Nos. 5,630,868; 5,559,169; 5,571,311; and 5,559,169 as well as U.S. patent application Ser. No. 08/572,525, all of which are hereby incorporated in their entireties by reference herein.\nThe particle having the attached vinyl substituted aromatic group or the vinyl substituted alkyl group is then reacted with a reactive free radical source and a stable free radical source for a sufficient time and at a sufficient temperature to form the modified particle having the attached group, like that of formula (I). The molar ratio of the reactive free radical source to the stable free radical source is from about 0.7 to about 2.0, and preferably from about 0.9 to about 1.1. Examples of the reactive free radical source include, but are not limited to, radicals generated from organic peroxides such as benzoyl peroxides and azo initiators such as azobisisobutyronitrile, and the like. The reactive free radical source is present in amounts sufficient to react with the vinyl aromatic group or vinyl alkyl group present on the particle. Preferably, the amount of the reactive radical source is from about 0.01 mmoles/g (of particle) to about 10 mmoles/g and more preferably from about 0.01 to about 5 mmoles/g.\nAlternatively, the modified particles of the present invention can be made by first forming the groups of the formulas described above, such as formula (I). Preferably, the stable free radical group is attached in the meta or para position of the aromatic group, if one is used. The group having the aromatic group or alkyl group with the stable free radical group is then attached to the particle by a diazonium treatment in the manner described in the above referenced patents and patent applications, where a diazonium salt will first be formed containing a group of one of the formulas described above in the manner described in the above-referenced patents. The groups of the formulas can be then subsequently attached to the particle. In a less preferred way, the formulas of the present invention can be attached to the particle through a hydroxyl or carbonyl group present on the particle. Also, the modified particle of the present invention can be formed by attaching a stable free radical compound containing at least one alkoxy silyl, alkoxy titanyl, or alkoxy zirconyl group to the particle which, in this particular process is preferably a metal oxide or metal hydroxide, or a carbon surface. This particular embodiment would attach a group of formula (II) or (IV) to a particle.\nIn another process, the modified particle can be formed by first taking an aromatic group or an alkyl group and reacting it with a compound to form the groups of the formulas described above except for the presence of the xe2x80x94SFR group. In other words, a group having the formula Axe2x80x94R1xe2x80x94CR2R3 would first be formed and then the xe2x80x94SFR would be attached to this group to form a group of formula (I) of the present invention which can then be attached to the particle. In this process, R2 and R3 are preferably hydrogen.\nA related process involves first taking an aromatic group or an alkyl group, such as a group having the formula Axe2x80x94R1xe2x80x94CR2R3 and attaching it onto the particle, for instance by a diazonium treatment in a manner described above, and then attaching the xe2x80x94SFR group to form a group of formula (I) of the present invention. Example 7 of the present application shows a specific embodiment of this process. In more detail, this process involves proton abstraction wherein a group containing an abstractable proton which leaves behind a carbon-centered radical is attached onto an aromatic group or an alkyl group, wherein the aromatic group or the alkyl group is directly attached to the particle. For purposes of this process, the aromatic group or the alkyl group can first be attached to the particle and then subsequently a group containing the abstractable proton can then be attached or an aromatic group or an alkyl group containing an abstractable proton can be attached onto the particle, for instance, by a diazonium treatment as described above. Then, a reactive free radical source and a stable free radical source are reacted with the particle having the attached aromatic group or alkyl group with the group containing an abstractable proton for a sufficient time and at a sufficient temperature to form the modified particle having the attached group, like that of formula (I). This modified particle having the attached group, like that of formula, (I) can then be subjected to polymerization reactions in order to have polymers attached thereto such as set forth in formula (III). In this process, the reactive free radical source abstracts the proton from the group containing the abstractable proton which thus leaves a carbon-centered radical behind which permits the stable free radical source to attached thereto to create the modified particles of the present invention such as set forth in formula (I). Examples of a group containing an abstractable proton which leaves behind a carbon-centered radical include alkyl groups, alkenes with a hydrogen atom in the alpha position to the olefinic bond, and the like. Specific examples include C1-C20 alkyl groups, such as methyl, ethyl, propyl, butyl, and the like. The reactive free radical source for this process can be organic peroxides such as benzoyl peroxides and t-butyl peroxides. Examples of the stable free radical source are described above and can be the same for this process. This process can also be applied where the modified particle aggregate has a group attached like that of formula (II) and (IV).\nAlternatively, the process of making the modified particles of the present invention can be a three step process where the alkyl or aromatic group is first attached to the particle and then the group having the formula R1xe2x80x94Cxe2x80x94R2R3 can be attached to the aromatic group or alkyl group. Then in a third step, the xe2x80x94SFR can be attached to the R1xe2x80x94Cxe2x80x94R2R3 group.\nAlso, in making the modified particles having the attached group of one the formulas described above, an aliphatic vinyl group can be attached to the particle surface by a diazotisation of a vinyl substituted amino aromatic or alkyl compound, or a vinyl substituted alkoxy silyl, titanyl, or zirconyl derivative. The vinyl group is then reacted with an organic peroxide and a stable free radical such that the organic peroxide and stable free radical are present in an amount sufficient to react with at least one percent of the vinyl groups and preferably from about 50 to about 100% of the vinyl groups and the mole ratios of the organic peroxide to stable free radical are preferably from about 0.5:1 to about 1:1.\nThe modified particles of formula (II) can be made in the following manner. The aggregate comprising at least one carbon phase and at least one metal-containing species phase can be made as described in U.S. patent application Ser. No. 08/828,725, hereby incorporated herein by reference. The aggregate comprising at least one carbon phase and at least one metal containing species phase can be made as described in U.S. patent application Ser. No. 08/446,141 and 08/750,017. The aggregate or particle can then be reacted with a coupling agent by adding the coupling agent to the aggregate in a medium and mixing. Then, the aggregate or particle having the attached coupling group can be reacted with a reactive radical source and stable free radical source as described above.\nFor purposes of the above-described reactions, the reaction should occur for a time and temperature sufficient to form the attached group onto the particle or aggregate. Generally, this time is from about 3 minutes to about 96 hours and more preferably from about I hour to about 24 hours. The temperature of the reaction is dependent on the half-life of the peroxide, preferably from about 50xc2x0 C. to about 200xc2x0 C., and more preferably from about 75xc2x0 C. to about 125xc2x0 C.\nFor the various methods described above, any solvent, aqueous or non-aqueous, can be used. Preferably, the solvent does not interfere with the radical on the particle. Preferably, the solvent is toluene, benzene, or xylene. Mixtures of various solvents can be used as well.\nWith the modified particles described above, polymers can be attached onto these attached groups by reacting these modified particles or aggregates with one or more polymerizable monomer(s) such as a vinyl or diene containing monomer. Specific examples of Such monomers include, but are not limited to styrene, isoprene, butadiene, chloromethylstyrene, methyl methacrylate, and butyl methacrylate, as well as acrylic acid and esters of acrylic acid and methacrylic acid and esters of methacrylic acid. More specific monomer examples include, but are not limited to, hydroxyalkyl (meth)acrylates, such as hydroxyethyl (meth)acrylate, hydroxyproyl (meth)acrylate; alkyl (meth)acrylates such as methyl acrylate or butyl acrylate; glycidyl (meth)acrylate; dimethylaminoethylacrylate; 2-acryl trimethyl ethylammonium chloride; aminoethyl acrylate; acrylonitrile; vinyl acetate, and the like. Mixtures of two or more monomers can be also used and/or polymerized sequentially. For purposes of the present invention, the vinyl containing monomer includes vinyl containing monomers having additional olefinic groups which are conjugated or unconjugated to the vinyl containing monomer.\nThe polymerization reaction is conducted under conditions which permits the polymerization of the monomer so that it forms a part of the group attached onto the particle or aggregate. These conditions are preferably heating modified particles with a monomer above 80xc2x0 C., preferably from about 120xc2x0 C. to about 150xc2x0 C., optionally in the presence of a solvent. The reaction can be ended by lowering the temperature below 80xc2x0 C. The polymer-modified particle can then be subjected to distillation, steam stripping, or precipitation or other known methods in order to remove residual monomers and solvents.\nThe polymerization reaction thus can produce any length of polymer on the modified particle or aggregate. For example, polymers having average molecular weights, preferably ranging from about 500 to about 1,000,000 can be made. Other ranges of average molecular weights include, but are not limited to, from about 500 to about 3,000, and from about 500 to about 2,000, and from about 500 to about 500,000, and from about 1,000 to about 250,000. The polymers can be any type, such as homopolymers, co-polymers, ter-polymers, or higher chain polymers. The polymers can also be block, graft, or random-type polymers. The polymers can be branched or unbranched. Polymers, such as comb polymers, can be formed or located off of the main polymer chain, and these polymers can be any type of polymer. Specific examples include, but are limited to, polyamines, such as polyethyleneimine or polyalkylene oxides, grafted onto the main polymer chain. The polymers that are formed can be ionic (e.g., cationic or anionic) or non-ionic polymers. Specific examples of polymers include, but are limited to, polyamides, polyacrylics, polyvinylacetates, polyvinylalcohols, polyvinylpyridines, polyvinylpyrrilodones, polyvinylinidazoles, and acrylainides and derivatives thereof as well as combinations thereof.\nOnce the polymerization occurs, the modified particle will have a group attached having the formula: \nwherein the substituents are the same as described earlier for formula (I) and X represents a polymer formed from at least one polymerizable vinyl or diene containing monomer.\nSimilarly, when the modified particle or aggregate having a group of formula (II) is polymerized by the introduction of one or more monomers, the group attached to the particle or aggregate will have the formula: \nwherein the substituents are the same as described in formula (II) and X represents a polymer formed from at least one polymerizable vinyl or diene containing monomer.\nIn another embodiment, a modified particle has an organic group containing a xe2x80x94SFR group directly attached to the particle. Preferably, the xe2x80x94SFR group is directly attached to the particle. The xe2x80x94SFR group and particle can be the same as discussed above. This type of modified particle can be prepared by heating a particle, like carbon black, with a organic group comprising a stable free radical in a solvent, such as toluene and preferably, in an inert atmosphere. With this modified particle having an attached organic group having a xe2x80x94SFR group, polymers can be attached onto the xe2x80x94SFR group by reacting the modified particle with a polymerizable monomer as discussed above using the same procedure.\nWith any of the above-described groups containing the xe2x80x94SFR group, the modified particle or aggregate containing the xe2x80x94SFR group or the polymerized versions thereof can be terminated by any means known to those skilled in the art in view of the present application so that a terminated moiety (-T) is located in the position of the xe2x80x94SFR group in the formulas above. In particular, the xe2x80x94SFR group can be replaced with a proton (e.g., hydrogen atom), subjected to disproportionation, or replaced with a chemical group through group transfer and the like. For instance, the xe2x80x94SFR group can be replaced with a metal containing group (e.g., tin or tin compound), a hydroxy group, or a halide group. For purposes of the present invention, the terminated moiety includes any chemical group capable of forming a chemical bond by substitution or replacement of the xe2x80x94SFR group, by one or multiple steps. For instance, the hydroxy group or halide group replacing the xe2x80x94SFR group can then be subsequently modified or replaced with chemical groups, such as organic groups, by such means as conventional oxidation reactions known to those skilled in the art. This termination of removing the xe2x80x94SFR group and replacing it or terminating it with another group can be accomplished, for instance, by methods described in 214 ASC Nat\"\"l Mtg, Las Vegas, Sep. 7-11, 1997, Paper in ORGN O61, Gravert, D. J.; Datta A.; Wentworth P., Jr.; Janda, K. D., which are incorporated in their entirety by reference herein.\nThe modified particles of the present invention, and preferably the polymerized and terminated versions of the modified particles of the present invention, can form part of a polymeric composition and be present with other ingredients commonly used and found in polymeric compositions.\nThe modified particle of the present invention can be used in a variety of applications. For instance, it can be used in coating formulations or compositions, ink formulations or compositions, or toner formulations or compositions, such as printing inks and inkjet inks, toners, automobile coatings, and the like. Also, the modified particles can be used as reinforcers for compositions, such as polymeric compositions and can also serve as impact modifiers, or as agents used to increase compatibility of a polymeric composition.\nIn more detail, reinforcement of elastomeric compositions including tire, hose, profile extrusion, seals, gaskets, and vibration isolation units, as well as the specific reinforcement of a single elastomer phase in a multiphase elastomer blend composition; reinforcement of thermoplastic compositions such as polyolefines, styrenic, acrylics, polyesters and polyamides, and thermoplastic elastomers and thermoplastic polyolefines; reinforcement of thermoset compositions, e.g., acrylics; impact modification of thermoplastic compositions; impact modification of thermosets; highly dispersible masterbatch for pigmentation, reinforcement, and/or UV protection of thermoplastic compositions, coatings, thermoplastic elastomers, and crosslinked compositions; as a synthetic support for solid phase organic synthesis; as a support or medium for effluent extraction processes - both organic and inorganic components; as a catalyst support; and/or as a superadsorbant for either aqueous of hydrocarbon materials, e.,q., use in sanitary wear, growing medium for plants.\nThe modified particles of the present invention can be incorporated in and form a part of elastomeric compositions. Other conventional ingredients for elastomeric compositions can also be present, such as coupling agents and the like.\nCoupling agents, as used herein, include, but are not limited to, compounds that are capable of coupling fillers such as carbon black or silica to an elastomer having one functionality which will attach to the particle and another functionality which will attach to the elastomer. Coupling agents useful for coupling silica or carbon black to an elastomer, are expected to be useful with the silicon-treated carbon blacks. Useful coupling agents include, but are not limited to, silane coupling agents such as bis(3-triethoxysilylpropyl)tetrasulfane (Si-69), 3-thiocyanatopropyl-triethoxy silane (Si-264, from Degussa AG, Germany), xcex3-mercaptopropyl-trimethoxy silane (A189, from Union Carbide Corp., Danbury, Conn.); zirconate coupling agents, such as zirconiumin dineoalkanolatodi(3mercapto) propionato-O (NZ 66A, from Kenrich Petrochemicals, Inc., of Bayonne, N.J.); titanate coupling agents; nitro coupling agents such as N,Nxe2x80x2-bis(2-methyl-2-nitropropyl)-1,6-diaminohexane (Sumifine 1162, from Sumimoto Chemical Co., Japan); polyalkoxysiloxane (e.g. Zeruma from the Yokohama Rubber Co. Ltd., Japan) and mixtures of any of the foregoing. The coupling agents may be provided as a mixture with a suitable carrier, for example X50-S which is a mixture of Si-69 and N330 carbon black, available from Degussa AG.\nThe elastomeric compounds of the present invention may be prepared from the modified particles of the present invention by compounding with any elastomer including those useful for compounding a carbon black.\nAny suitable elastomer may be compounded with the modified particles to provide the elastomeric compounds of the present invention. Such elastomers include, but are not limited to, homo- or co-polymers of 1,3 butadiene, styrene, isoprene, isobutylene, 2,3-dimethyl-1,3-butadiene, acrylonitrile, ethylene, and propylene. Preferably, the elastomer has a glass transition temperature (Tg) as measured by differential scanning colorimetry (DSC) ranging from about xe2x88x92120xc2x0 C. to about 0xc2x0 C. Examples include, but are not limited, styrene-butadiene rubber (SBR), natural rubber, polybutadiene, polyisoprene, and their oil-extended derivatives. Blends of any of the foregoing may also be used.\nAmong the rubbers suitable for use with the present invention are natural rubber and its derivatives, such as chlorinated rubber and epoxidized rubber. The modified particles of the present invention may also be used with synthetic rubbers such as: copolymers of from about 10 to about 70 percent by weight of styrene and from about 90 to about 30 percent by weight of butadiene such as copolymer of 19 parts styrene and 81 parts butadiene, a copolymer of 30 parts styrene and 70 parts butadiene, a copolymer of 43 parts styrene and 57 parts butadiene and a copolymer of 50 parts styrene and 50 parts butadiene; polymers and copolymers of conjugated dienes such as polybutadiene, polyisoprene, polychloroprene, and the like, and copolymers of such conjugated dienes with an ethylenic group-containing monomer copolymerizable therewith such as styrene, methyl styrene, chloromethylstyrene, acrylonitrile, 2-vinyl-pyridine, 5-methyl 2-vinylpyridine, 5-ethyl-2-vinylpyridine, 2-methyl-5-vinylpyridine, alkyl-substituted acrylates or methacrylates, vinyl ketone, methyl isopropenyl ketone, methyl vinyl ether, alpha-methylene carboxylic acids and the esters and amides thereof such as acrylic acid, methacrylic acid, and dialkylacrylic acid amide; also suitable for use herein are copolymers of ethylene and other high alpha olefins such as propylene, butene, pentene, hexene, and octene. Other monomers that could be used include norbornene and hex-1,5-diene, and the like.\nThe rubber compositions of the present invention can therefore contain one or more elastomers, curing agents, reinforcing fillers, coupling agents, and, optionally, various processing aids, oil extenders, and antidegradents. In addition to the examples mentioned above, the elastomer can be, but is not limited to, polymers (e.g., homopolymers, copolymers, and terpolymers) manufactured from 1,3 butadiene, styrene, isoprene, isobutylene, 2,3-dimethyl-1,3-butadiene, acrylonitrile, ethylene, propylene, and the like. It is preferred that these elastomers have a glass transition point (Tg), as measured by DSC, between xe2x88x92120xc2x0 C. and 0xc2x0 C. Examples of such elastomers include poly(butadiene), poly(styrene-co-butadiene), and poly(isoprene).\nElastomeric compositions disclosed in the present invention include, but are not limited to, vulcanized compositions (VR), thermoplastic vulcanizates (TPV), thermoplastic elastomers (TPE) and thermoplastic polyolefins (TPO). TPV, TPE, and TPO materials are further classified by their ability to be extruded and molded several times without loss of performance characteristics.\nThe elastomeric compositions may include one or more curing agents such as, for example, sulfur, sulfur donors, activators, accelerators, peroxides, and other systems used to effect vulcanization of the elastomer composition.\nThe resultant elastomeric compounds containing the aggregates of the present invention and optionally containing one or more coupling agents may be used for various elastomeric products such as a tread compound, undertread compound, sidewall compound, wire skim compound, innerliner compound, bead, apex, any compound used in carcass and other components for vehicle tires, industrial rubber products, seals, timing belts, power transmission belting, and other rubber goods.\nIn an embodiment of the present invention, the elastomeric compositions of the present invention, which contain at least one modified particle of the present invention, can have a total residue after nitrogen pyrolysis at 650xc2x0 C. of from about 1% to about 60% by weight.\nFor purposes of the present invention, the ink includes at least one type of polymerized-terminated modified particle, and at least one ink vehicle. The inkjet ink formulation includes at least one type of polymerized terminated modified particle, and at least one inkjet ink vehicle. Generally, the formulations described in WO 97/47699 with respect to coatings and inks can be used herein, however, incorporating the modified particles of the present invention. The coating formulations contain at least one type of polymerized-terminated modified particle and at least one suitable solvent. Other conventional additives may be incorporated into the coating formulations, such as a binder.\nEach of these formulations can contain additional conventional colorants and other optional, conventional ingredients or additives, for instance as described in U.S. Pat. Nos. 5,571,311; 5,672,198; 5,266,361; 5,707,432 such as a humectant, binders, dyes, biocides, surfacants, penetrants. All of these patents are incorporated herein in their entirety by reference herein.\nThe toner can contain at least one polymerized-terminated modified particle and resin particles.\nThe conventional and standard ingredients or additives for toner formulations, such as those described in U.S. Pat. Nos. 5,278,018; 5,275,900; 5,695,899 and 5,116,712, can be used in the present invention, and are incorporated in their entirety by reference herein. The toner can be a positively or negatively charged toner composition.\nIn the various products incorporating the modified particles of the present invention, one or more types of other particles can also be present, such as a carbon product (e.g., carbon black), metal oxide (e.g., silica, zinc oxide, and aluminum oxide), metal silicate (e.g., clay, aluminum, calcium, magnesium silicates), modified silica, modified carbon products having an attached organic group, an aggregate comprising at least one carbon phase and at least one silicon-containing species phase optionally having an attached organic group, carbon black at least partially coated with silica, or any combination thereof. In any product mentioned herein, more than one type of modified particle of the present invention can be present in the various products described above."}
{"text": "Many types of electrical equipment contain IC (integrated circuit) devices which are vulnerable to damage from high voltage transients.\nIn a television receiver the anode of the image producing kinescope is typically biased at a high potential, e.g., 25,000 volts. High voltage transients may be produced when the high voltage anode of the kinescope is rapidly discharged to points at lower potentials. Such high voltage transients have positive and negative peaks often in excess of 100 volts and may last several microseconds. High voltage transients may also be produced when electrostatic charges are discharged as a user contacts the controls of the television receiver. High voltage transients may be coupled to the terminals of IC's employed in the television receiver for video and audio signal processing. Accordingly, these IC's may be damaged by high voltage transients.\nCopending patent application entitled \"Protection Circuit for Integrated Circuit Devices,\" Ser. No. 212,534, filed in the name of the present inventor Dec. 3, 1980 and assigned to RCA Corporation, discloses a protection device for protecting an IC from negative voltage transients that exceed the most negative power supply potential applied to the IC. Another copending patent application entitled \"Protection Circuit for Integrated Circuit Devices,\" Ser. No. 230,357 filed also in the name of the present inventor Jan. 30, 1981, and also assigned to RCA Corporation, discloses a protection circuit which is triggered at one forward biased diode voltage drop above the positive supply voltage for protecting an IC from positive voltage transients that exceed the most positive power supply potential applied to the IC.\nHowever, in a television receiver, particular signals applied to an IC may have positive voltage excursions which in normal operation exceed the positive supply potential. For example, a typical television horizontal/vertical regulator IC requires a feedback connection from the kinescope deflection coils to one of its input terminals. While the power supply for the IC is typically +10 volts, the peak feedback voltage from the deflection coils is typically +27 volts. Therefore, it is desirable to provide a positive transient protection circuit for such IC's that permits normal signal voltages to exceed the power supply potential without activating such protection circuit and nevertheless protects the IC from excessively large transients."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of electronic design automation systems for designing and characterizing integrated circuits. More specifically, the present invention relates to an evaluation system for evaluating a technology library for transformation into an effective modeling method and data structure for modeling signal propagation delay of a timing arc of an integrated circuit cell.\n2. Related Art\nThe rapid growth of the complexity of modern electronic circuits has forced electronic circuit designers to rely upon computer programs to assist and automate most steps of the circuit design process. Typical circuits today contain hundreds of thousands or millions of individual pieces or xe2x80x9ccells.xe2x80x9d Such a design is much too large for a circuit designer or even an engineering team of designers to manage effectively manually. To automate the circuit design and fabrication of integrated circuit devices, electronic design automation (EDA) systems have been developed.\nAn EDA system is a computer software system designers use for designing integrated circuit (IC) devices. The EDA system typically receives one or more high level behavioral descriptions of an IC device (e.g., in HDL languages like VHDL, Verilog, etc.) and translates this high level design language description into netlists of various levels of abstraction. At a higher level of abstraction, a generic netlist is typically produced based on technology independent primitives. The generic netlist can be translated by the EDA system into a lower level technology-specific netlist based on a technology-specific library that has gate-specific models for timing and power estimation. A netlist describes the IC design and is composed of nodes (elements) and edges, e.g., connections between nodes, and can be represented using a directed cyclic graph structure having nodes which are connected to each other with signal lines. The netlist is typically stored in computer readable media within the EDA system and processed and verified using many well known techniques. One result is a physical device layout in mask form which can be used to directly implement structures in silicon to realize the physical IC device.\nMore specifically, within a typical EDA system, the circuit designer first produces a high-level description of the circuit in a hardware description language such as Verilog or VHDL. This high-level description is converted into a netlist using a computer implemented synthesis process such as the xe2x80x9cDesign Compilerxe2x80x9d by Synopsys of Mountain View, Calif. A netlist is a description of the electronic circuit which specifies what cells compose the circuit and which pins of which cells are to be connected together using wires (xe2x80x9cnetsxe2x80x9d). Importantly, the netlist does not specify where on a circuit board or silicon chip the cells are placed or where the wires run which connect them together. Determining this geometric information is the function of a computer controlled placement process and a computer controlled routing process.\nThe placement process finds a location for each cell on a circuit board or silicon chip. The locations are specified, typically, in two dimensional spatial coordinates, e.g., (x, y) coordinates, on the circuit board or silicon chip. The locations are typically selected to optimize certain objectives such as wire length, wire routibility, circuit speed, circuit power consumption, and/or other criteria, subject to the condition that the cells are spread evenly over the circuit board or silicon chip and that the cells do not overlap with each other. The output of the automatic computer controlled cell placement process includes a data structure including the (x, y) location for each cell of the IC design.\nNext, the designer supplies the netlist and the cell location data structure, generated by the placement program, to a computer implemented automatic wire routing process (xe2x80x9crouterxe2x80x9d). The router generates wire geometry within data structure for connecting pins together. The wire geometry data structure and cell placement data structure together are used to make the final geometric database needed for fabrication of the circuit. Routers typically include a coarse routing process and a fine routing process. The coarse router provides a general path for the routing that is done at the detail stage. The coarse router examines at the level of the whole integrated circuit chip and its available resources and determines what the rough pathways should be from a topological standpoint. The fine or detail router lays down the actual geometries and connected wire segments in the appropriate layers as a wire connection may span multiple layers. The fine router creates wire routes that are xe2x80x9cclean,xe2x80x9d e.g., do not have design rule violations, do not overlap other structures and can be fabricated.\nThe signal propagation delay (xe2x80x9ccell delayxe2x80x9d) through a cell (xe2x80x9cgatexe2x80x9d) is an important characteristic to model within an EDA system. The cell delays in a technology library are typically represented using non-linear delay models (NLDM) which are essentially look-up tables. Typically, a group of tables are supplied in the technology library for each cell, tables are designated for representing the rise and fall delays for each timing arc of the cell. These tables are typically 4-dimensional in that they accept output load and input transition time (slew) as inputs and generate delay and output slew values as outputs. These output load-based NLDMs, while providing delay values, have a disadvantage in the cell delay modeling processes that are performed early in the circuit synthesis process. For instance, during early circuit synthesis processes, the output load of the cell is not known because the cells have not yet been mapped to the target technology library and, as such, the cells are not yet connected together. Output load estimates are made in these early synthesis processes because the output load-based NLDMs need these values as inputs. Unfortunately, the output load estimates introduce inaccuracies in the overall circuit synthesis process. The output load-based NLDMs also introduce a xe2x80x9cCatch-22xe2x80x9d problem in that delay modeling helps to accurately map the cells, but mapping yields the true output capacitance that is then used to accurately determine the cell\"\"s delay, etc. It would be advantageous to provide a cell delay modeling system that did not require the output load of a cell as an input.\nA more simplistic linear delay model has been proposed as a vehicle for efficient logic synthesis of high-performance designs. This delay model is also referred to as the constant delay model. In the constant delay model, the delay of a timing arc of a gate, xcfx84, is represented as:\nxcfx84=Rxe2x80xa2Co+p\nwhere R=output resistance,\nCo=output load of the timing arc output, and\np=intrinsic delay of a gate.\nThis relationship can also be represented as:\nxcfx84=(Rxe2x80xa2Ci)xe2x80xa2(Co/Ci)+p\nwhere Ci=input capacitance of the timing arc input.\nThe term (Rxe2x80xa2Ci) is also referred to as the logical effort of the gate. The term (Co/Ci) is also referred to as the electrical effort or gain of the gate.\nThe constant delay model assumes that the delay of the timing arc remains constant. The reasoning for this is as follows. The intrinsic delay of the gate, p, is constant. As Co increases, the gate is implicitly upsized, so Ci increases appropriately. So, (Co/Ci) remains constant. As Ci increases, R appropriately decreases, so (Rxe2x80xa2Ci) remains constant. Consequently, xcfx84 remains constant. An important property of the constant delay model is that the delay of a timing arc is independent of load, e.g., the delay does not depend on either Co or Ci, but merely on the ratio of the two. This property is useful in early stages of logic synthesis, prior to technology dependent optimization because the actual load of a gate is unknown at that time.\nHowever, while this simplistic delay model does have its advantages, the constant delay model does not consider several important factors when modeling the delay of a gate. For instance, the constant delay model does not consider the impact of transition times on delay, nor does it deal effectively with complex gates with different input capacitances for different input pins. The model also fails to take into consideration the impact of limited discrete sizes in the technology library nor of certain design rules like maximum capacitance and maximum transition associated with gates. By not taking into consideration the above factors, the constant delay model is not as accurate as non-linear delay models. It would be advantageous to provide a method and system for providing cell delay modeling that offered the advantageous of constant cell delay modeling but also considered the above referenced factors.\nA number of technology libraries currently exist that utilize output load-based non-linear delay models (NLDMs). It would be advantageous to provide a mechanism that automatically constructed a technology library that includes output load-based NLDMs and could also be transformed into an alternate library format that used delay models that do not require output load as an input, thereby alleviating the above problems. If such a library construction process or system existed, vendors of technology libraries could readily construct technology libraries can be used with converters to the alternate library format.\nA system and computer implemented process are described herein for constructing a technology library that is suitable for use with an electronic design automation system that converts the target technology library into a scalable cell library having non-linear, gain-based delay models for estimating circuit delay. The scalable cell library can then be used by gain-based structuring and mapping processes. In one embodiment, the library construction process places at least six discrete cells in each logic function of a basic cell set. The library construction process also places at least five discrete cells in each logic function of an extended cell set and rules out cell sizing using buffer circuits. Also, for each discrete cell in the complete cell set, the variance of the capacitances between different input pins of the cell is maintained to be within 10 percent. For corresponding timing arcs of discrete sizes for a particular logic function, the present invention keeps the delay of the corresponding timing arcs within 10 percent of each other, when the ratio of the output load to input capacitance are equal for the corresponding input-output pin paris. Also, the present invention constructs a technology library that has geometrically distributed sizes of cells within each logic function. Lastly, for each discrete cell within a logic cluster, the output maximum capacitance constraint is kept linearly proportional to the average input capacitance of the discrete cell. These processes likely allow a technology library to be suitable for the generation of a scalable library which can be used for integrated circuit design and fabrications.\nEmbodiments of the present invention also include a library evaluation system for determining the suitability of a target technology library for generating a suitable scalable library that utilizes delay models that are gain-based. Other embodiments of the present invention provides a cell delay modeling system and method for modeling cell delay without requiring the output load capacitance of the cell as an input. Furthermore, the present invention provides a non-linear delay model for accurate cell delay modeling. The present invention provides a cell delay modeling system and method that considers the impact of transition times on delay, that deals effectively with complex gates with different input capacitances for different input pins, that takes into consideration the impact of limited discrete sizes in the technology library and furthermore takes into consideration certain design rules like maximum capacitance and maximum transition associated with gates. These and other advantages of the present invention not specifically mentioned above will become clear within discussions of the present invention presented herein.\nAn evaluation system is described for evaluating the suitability of a target technology library for use with an electronic design automation system that converts the target technology library into a scalable library having non-linear, gain-based delay models for estimating circuit delay. After a library analysis process executes, an internal scalable library is generated having a scalable cell model for each functional cell class of the target technology library. Internal characteristics of the scalable library are analyzed to determine whether or not the target technology library is suitable for generating a suitable scalable library. Suitable scalable libraries are used for gain-based structuring and mapping processes. The library evaluation processes uses several metrics to determine the suitability of the target technology library. The first metric is the number of sizes metric and the second metric is the size consistency metric. The number of sizes metric is computed with respect to subsets (e.g., basic set and extended set) of the scalable cells identified within the scalable library. The number of sizes metric determines whether or not there are sufficient numbers of discrete cells within the cell clusters of the subsets. The size consistency metric examines the input pin capacitance ratios and the delay error amounts for discrete cells of the scalable library. An evaluation report is then compiled using the above metric information to grade library suitability.\nIn another embodiment, a non-linear, gain-based modeling of circuit delay within an electronic design automation (EDA) system is described herein. The present invention provides a scalable cell model for use in logic structuring and mapping for the design of integrated circuits. The scalable cell model provides a four dimensional non-linear delay model that accepts (1) input slew and (2) gain and provides (3) delay and (4) output slew information for a scalable cell. By eliminating output loading as a requirement for delay computations, the scalable model of the present invention can effectively be used to provide accurate delay (timing) information for those logic synthesis processes that precede technology based optimizations where the actual load of a cell is unknown. This scalable cell model considers: the impact of transition times on delay; complex gates having different input capacitances for different input pins; the impact of limited discrete cell sizes in the technology library; and design rules, e.g., maximum capacitance and maximum transition associated with gates.\nA library analysis process creates a scalable technology library by analyzing discrete cells from a target technology library. The target technology library is analyzed and clustering is performed to select a cluster (e.g., subset) of cells for each cell group sharing a common functionality. Clustering is performed based on a cluster metric that can be any of the following: a delay consistency metric, an input pin consistency metric, an intercept/slope consistency metric or an consistency metric based on the slope and average input capacitance being inversely proportional. A nominal input slew value is computed for all cells and a scaling factor is also computed for each cell of each cell cluster. From each cell cluster, a four dimensional scalable cell model (look-up table) is then generated that inputs gain and input slew and outputs delay and output slew. Separate non-linear delay models are provided for rise and fall time for each cell. Those cells not within a cell cluster do not participate in generating data for the scalable model. Therefore, the scalable cells in the scalable library are characterized with a new load-independent delay model, in which delay is modeled as a non-linear function of gain and slew.\nA default gain is then computed for each scalable cell model and an area model and an input pin capacitance model are then generated for each scalable cell model. The input pin capacitance model is useful for load propagation of the scalable cells. For a particular integrated circuit design, its designated discrete cells are replaced with the scalable cells which are used during logic structuring and mapping. Logic structuring and mapping processes use the scalable cells as a basis for the optimization. After synthesis, a discretization process then converts the scalable cells back to discrete cells before the technology dependent optimizations are performed. Therefore, at the end of optimization, present invention discretizes the scalable cells to the closest discrete cells available in the target technology library. Following this, technology based optimization is performed primarily for delay, design rule fixing and area recovery. The accuracy of the logic structuring and mapping processes is hinged upon the accuracy of the scalable model that library analysis derives.\nThe present invention also has a capability of automatically determining whether a technology library is suitable for its optimization techniques. This is important because the accuracy of many processes in the present invention may depend on the accuracy of the gain-based scalable delay model that is created. After performing library analysis, a library evaluation step is performed to determine if the target technology library is suitable. If the library is found to be suitable, logic structuring and mapping is invoked and if not, a conventional optimization engine can be used."}
{"text": "1. Technical Field\nThe invention relates to reinforced resilient pneumatic tires and more particularly to a vehicle tire reinforced by a thin annular composite band which is stabilized by a plurality of radial elements in the tire sidewalls to enable the tire to run in an unpressurized condition. More particularly, the invention relates to a pneumatic tire in which the band element is formed by a plurality of helically wound layers of material strips, in which at least the outer strips are prestressed in tension to improve the endurance of the band and minimize band stresses placed thereon during both pressurized and unpressurized conditions.\n2. Background Information\nVarious tire constructions have been devised over the years which enable a tire to run in an under-inflated or non-inflated condition such as after receiving a puncture and loss of pressurized air for extended periods of time and at relatively high speeds. This enables the vehicle operator to safely drive the vehicle to an appropriate location for repair or replacement of the punctured tire. Certain of these safety tires, referred to as \"run flat tires\", have been successful for certain applications and certain types of tire constructions. Most of these run flat tires achieve their run flat capability, by the placement of reinforcing layers or members of relatively stiff elastomeric material in the side walls of the tire which enable the tire to support the vehicle weight even with the complete loss of internal air pressure. Examples of such prior art run flat tire constructions which use such sidewall inserts are shown in U.S. Pat. Nos. 3,911,987; 3,949,798; 3,954,131; 4,067,372; 4,202,393; 4,203,481; 4,261,405; 4,265,288; 4,287,924; 4,365,659; 4,917,164; and 4,929,684.\nIn addition to these prior art run flat tire patents, various run flat tire constructions have been developed which utilize a thin annular band which extends circumferentially throughout the tire beneath the tread area. Examples of such banded run flat tires are shown in the following patents.\nU.S. Pat. No. 4,428,411 describes a method to make a particular band for use in a run flat tire which uses a series of side-by-side elements in the form a helix. The band has hoop compression as against a conventional breaker belt that has no significant compressive strength but is used only to resist tension loads endured by the tire when pressurized. However, the band has no residual stresses after manufacture which reduces the critical operational stresses after integration into a tire as is the band of the present invention.\nU.S. Pat. Nos. 4,673,014 and 4,794,966 teach a method to acquire desirable prestressing in a fabricated band made of helical elements. Physically bending the larger diameter helix element around a smaller mandrel and securing it with a resin impregnated tape does acquire a desirable level of prestressing. However, the band of the present invention is achieved by inherent thermal shrinkage as the sequentially overlapped material layers eliminate any physical bending of a structural element as is required in the method of these two patents.\nU.S. Pat. No. 4,456,048 teaches a method of acquiring a change in band stiffness as a function of deflection. The band has a lower stiffness for normal pressurized operation and has a higher stiffness to support load when the tire is uninflated and experiences larger deflection. The band of the present invention does not have a variable stiffness but is prestressed and does not change the band stiffness. The band does favorably affect the range of operating stresses. Japanese Patent application No. JP 63141809 discloses a run flat tire having a banded element which is formed of layered strips of materials, such as an aramide filament which is impregnated with a high elasticity epoxy resin, which after hardening provides a stiffened band but provides no suggestion for prestressing a band in tension which is achieved during the manufacture of the band of the present invention. Thus the tire of this disclosure requires that elastomeric side wall inserts be utilized in combination with the band in order to achieve the desired run flat characteristics.\nOther run flat banded pneumatic tires are shown in U.S. Pat. Nos. 4,111,249; 4,318,434; 4,428,411; 4,459,167; and 4,734,144.\nBanded tires have been fabricated with band materials made of steel, aluminum, titanium, and epoxy and thermoplastic composites with glass, KEVLAR (aromatic polyamide) and graphite fiber reinforcement. The common failure mode with the light, economical laminate band construction is interlaminar shear within the band's primary bending neutral axis. This is a fatigue failure and is directly related to the spectrum of cyclic operating stress. As in all fatigue failures, the lower the stress, the longer the operating life.\nPrestressing techniques lower the critical operating stress. This critical operating stress occurs in the footprint of the tire and is induced by road surface irregularities, aggravated by ride dynamics. This feature was identified by operating loaded tires on dynometer wheels and then compared with actual road tests. Dynometer tests showed an increase in cyclic life by a factor of 10-20 times. Thus, anomalies in the road surface can greatly shorten tire life when compared to dynometer testing.\nThere are three recognized approaches to improve the interlaminar fatigue life of the bands: lower the critical operating stress, improve the toughness features of the matrix and the introduction of reinforcing fibers thru the failure plane. This latter approach uses a weaving process and is applicable to wet or preimpregnated fiber tows or thermoplastic tape.\nThe proposed approach of the present invention uses a prestressing technique to reduce the critical operating stresses. The unique feature identified in this approach is that of beneficial residual stresses are introduced in the band by the fabrication process. The sequential application of reinforced thermoplastic tapes employing heat and high pressure to the contacting surfaces results in a gradual buildup of residual stress. The geometric features of the invention retain the desired strength of the reinforcing fiber in the hoop direction and provide adequate lateral strength, with the level of beneficial residual stress being achieved in the preferred embodiment with a winding angle preferably of .+-.15.degree.-20.degree. off circumferential."}
{"text": "This invention relates generally to integrated circuit devices made up of more than one die.\nIn a number of integrated circuit applications, relatively complex functions may be involved. Thus, in some cases, the capabilities desired to be integrated into a single integrated circuit die may exceed the available processing capabilities. Some functionalities are placed on one die and other functionalities are placed on another die.\nThe packaged dies may then be coupled together, for example by securing them to printed circuit boards having metallic interconnections to connect signals travelling between the two dies. This has the disadvantage that the footprint or size of the combined integrated circuit device is increased because spacing is needed between the dies to allow the desired interconnections. In addition, the resulting device may be slower due to the impedances arising from the metallic interconnections and the contacts to each die.\nThus, there is a continuing need for techniques for providing advanced functions in compact, multiple die integrated circuits."}
{"text": "A common construction of a header for an agricultural implement includes either a cutter bar with a pickup reel and/or a rotating tine pickup belt for feeding material to an auger which collects the material into a rear discharge of the header. In certain types of crops, material can be thrown upward onto the frame of the header above the rear discharge opening. The material can then obstruct the view of the operator of the implement, or become entangled in the implement, or be wasted by being thrown onto the ground instead of being collected in the usual manner of the implement."}
{"text": "A conventional rotation angle detection apparatus for detecting a rotation angle of a detected object includes a magnetism generation section, such as a magnet, and a magnetic flux density detection section, such as a Hall element. One of the magnetism generation section and the magnetic flux density detection section is attached to the detected object. The magnetic flux density detection section detects a magnetism of the magnetism generation section when the detected object rotates, thereby detecting the rotation angle of the detected object. For example, JP-A-2002-206911 discloses a rotation angle detection apparatus that includes a magnet as a magnetism generation section, a magnetic flux density sensor as a magnetic flux density detection section, and a signal processor as a processing section. The signal processor executes a correction operation. The signal processor corrects actual output voltages, which are based on an output signals of the magnetic flux density sensor, on the basis of predetermined values (voltage level) corresponding to predetermined correction points. The predetermined correction points are set at regular intervals.\nIn cases where the intervals of the correction points are constant, a nonlinearity of the output signal is not improved, and an error in the actual output voltages remains as a linearity error in the rotation angle. JP-A-2007-71889 discloses a correction method that divides a region where an output change rate is large more finely than a region where the output change rate is small. Japanese Patent No. 3,491,577 (corresponding to U.S. Pat. No. 6,498,479) discloses a correction method in which output signals of a magnetism detection element are converted into digital signals, and the digital signal are processed with an arcsine function. However, when an output waveform is unclear, the correction methods disclosed in JP-A-2007-71889 and Japanese Patent No. 3,491,577 cannot be used."}
{"text": "Not Applicable\nNot Applicable\nThe present invention relates to fan drive assemblies of the type including a cooling fan and a fan drive, and more particularly, to such fan drive assemblies wherein the fan drive is of the type in which heat is generated as a result of the transmission of torque within the fan drive, and the ability of the fan drive to dissipate such generated heat represents a limiting factor on the torque transmitting capability of the fan drive assembly.\nAlthough the present invention may be used with various types and configurations of torque transmitting fan drives, it is especially adapted for use with fan drive assemblies of the type including a viscous fluid coupling device as the fan drive, and will be described in connection therewith.\nFan drive assemblies of the type which may benefit from the use of the present invention have found several uses, one of the most common of which is in connection with cooling the radiator of a vehicle engine. As is well known to those skilled in the art, the fan drive of the typical fan drive assembly comprises a viscous fluid coupling device, so named because the coupling utilizes a high viscosity fluid to transmit torque, by means of viscous shear drag, from an input coupling member (clutch) to an output coupling member (housing), with the cooling fan being bolted, or otherwise suitably attached, to the output coupling member.\nThe present invention is especially advantageous when used on a relatively high horsepower fan drive assembly, i.e., one which is capable of transmitting somewhere in the range of about two to about twelve horsepower from the fan drive to the cooling fan. Typically, such high horsepower fan drives include an output coupling member of the type which comprises a cast aluminum body and a cast aluminum cover. The input coupling member is typically also made as a cast aluminum member, and cooperates with the body and/or the cover to define a plurality of interdigitated lands and grooves which define the viscous shear space. When the shear space is filled with viscous fluid, typically a silicon fluid, torque is transmitted from the input coupling member to the output coupling assembly, in response to the rotation of the input coupling member.\nDuring such torque transmission, substantial heat is generated as a result of the shearing of the viscous fluid between the lands and grooves. The amount of heat generated is generally proportional to the xe2x80x9cslip speedxe2x80x9d of the fan drive, i.e., the difference between the speed of the input and the speed of the output. It is generally well understood by those skilled in the art that the ability to transmit torque is limited by the ability of the device to dissipate the heat generated. For example, in a viscous fan drive, if the temperature of the viscous fluid exceeds a certain maximum temperature, the result will be a deterioration of the viscous properties of the fluid, resulting in a gradual loss of the torque transmitting capability of the fluid.\nIn the fan drive art, it has been conventional for the design and development of a particular cooling fan to occur generally independently of the design and development of the viscous fan drive with which the fan is to be utilized. In other words, the fan is designed to provide the desired operating parameters (e.g., torque, air flow, etc.), thus determining the blade configuration and spacing, and then the mounting portion of the fan (the xe2x80x9cspiderxe2x80x9d) is designed or merely modified to adapt to the configuration of the particular fan drive mounting arrangement (e.g., mounting pads or bosses, disposed at a particular diameter from the axis of the fan drive).\nWhat has not been conventional in the fan drive art is to design the cooling fan and the fan drive as a xe2x80x9cpackagexe2x80x9d, with the goal of maximizing the heat dissipation of the overall fan drive assembly. As a result, it would appear that, at the time of the present invention, there is no commercially available fan drive assembly which achieves nearly its optimum, potential heat dissipation (heat rejection). As a further result, practically every fan drive assembly in commercial use is larger and more expensive than is actually necessary, in order to achieve a particular, desired flow of cooling air through the radiator.\nAlthough the present invention is not limited to a fan drive assembly in which the fan is mounted to the rearward side of the housing (body), rather than being mounted to the cover, the invention is especially advantageous in such an arrangement, and will be described in connection therewith. A typical xe2x80x9crear mountxe2x80x9d fan drive is illustrated and described in U.S. Pat. No. 4,384,824, in which the body member includes four mounting bosses located radially inward of the body cooling fins. As a result, the fan spider interferes with the radial flow of cooling air through the body cooling fins, thus reducing the heat dissipation capability of the particular fan drive assembly.\nAccordingly, it is an object of the present invention to provide an improved fan drive assembly in which the cooling fan and the fluid coupling device driving the fan are designed such that the overall assembly approaches the optimum, potential heat dissipation.\nIt is a more specific object of the present invention to provide an improved fan drive assembly in which the cooling fan is mounted to the body (housing) of the fluid coupling device in a manner which substantially improves the flow of air through the housing cooling fins.\nIt is a related object of the present invention to provide an improved fan drive assembly which accomplishes the above-stated objects, and in which the cooling fan is configured to further improve the flow of air through the housing cooling fins.\nIt is another object of the present invention to provide an improved fan drive assembly which accomplishes the above-stated objects, and in which the cover cooling fins are configured to improve the flow of air through the cover cooling fins.\nThe above and other objects of the invention are accomplished by the provision of a fan drive assembly of the type comprising a cooling fan attached to a fluid coupling device, the cooling fan comprising a fan hub, a spider portion, and a plurality of fan blades extending radially from the fan hub. The fluid coupling device comprises a first rotatable coupling assembly including a body member having a rearward surface, and a cover member cooperating with the body member to define a fluid chamber therebetween, a second rotatable coupling member being disposed in the fluid chamber for rotation relative to the first coupling assembly. The first coupling assembly and the second coupling member cooperate to define a viscous shear chamber therebetween, whereby torque may be transmitted from the second coupling member to the first a coupling assembly in response to the presence of viscous fluid in the shear chamber. The body member includes a plurality of cooling fins and a plurality of mounting portions, the spider portion being attached to the mounting portions and defining a pilot diameter.\nThe improved fan drive assembly is characterized by the body member including a plurality of mounting portions, each of which is disposed immediately adjacent an outer periphery of the body member. Each of the mounting portions defines a machining chucking surface, and a spider mounting surface on a rearward face thereof. The spider mounting surface includes a pilot surface in engagement with the pilot diameter of the spider portion. The plurality of cooling fins covers substantially all of the rearward surface of the body member not covered by the mounting portions.\nIn accordance with another aspect of the present invention, the fan hub and each of the plurality of fan blades cooperate to define a rearward axially extending air dam portion operable to restrict localized radial air flow."}
{"text": "In the field of microprocessors, the number of instructions executed per second is a primary performance measure. As is well known in the art, many factors in the design and manufacture of a microprocessor impact this measure. For example, the execution rate depends quite strongly on the clock frequency of the microprocessor. The frequency of the clock applied to a microprocessor is limited, however, by power dissipation concerns and by the switching characteristics of the transistors in the microprocessor.\nThe architecture of the microprocessor is also a significant factor in the execution rate of a microprocessor. For example, many modern microprocessors utilize a \"pipelined\" architecture to improve their execution rate if many of their instructions require multiple clock cycles for execution. According to conventional pipelining techniques, each microprocessor instruction is segmented into several stages, and separate circuitry is provided to perform each stage of the instruction. The execution rate of the microprocessor is thus increased by overlapping the execution of different stages of multiple instructions in each clock cycle. In this way, one multiple-cycle instruction may be completed in each clock cycle.\nBy way of further background, some microprocessor architectures are of the \"superscalar\" type, where multiple instructions are issued in each clock cycle for execution in parallel. Assuming no dependencies among instructions, the increase in instruction throughput is proportional to the degree of scalability.\nAnother known technique for improving the execution rate of a microprocessor and the system in which it is implemented is the use of a cache memory. Conventional cache memories are small high-speed memories that store program and data from memory locations which are likely to be accessed in performing later instructions, as determined by a selection algorithm. Since the cache memory can be accessed in a reduced number of clock cycles (often a single cycle) relative to main system memory, the effective execution rate of a microprocessor utilizing a cache is much improved over a non-cache system. Many cache memories are located on the same integrated circuit chip as the microprocessor itself, providing further performance improvement.\nAccording to each of these architecture-related performance improvement techniques, certain events may occur that slow the microprocessor performance. For example, in both the pipelined and the superscalar architectures, multiple instructions may require access to the same internal circuitry at the same time, in which case one of the instructions will have to wait (i.e., \"stall\") until the priority instruction is serviced by the circuitry.\nOne type of such a conflict often occurs where one instruction requests a write to memory (including cache) at the same time that another instruction requests a read from the memory. If the instructions are serviced in a \"first-come-first-served\" basis, the later-arriving instruction will have to wait for the completion of a prior instruction until it is granted memory access. These and other stalls are, of course, detrimental to microprocessor performance.\nIt has been discovered that, for most instruction sequences (i.e., programs), reads from memory or cache are generally more time-critical than writes to memory or cache, especially where a large number of general-purpose registers are provided in the microprocessor architecture. This is because the instructions and input data are necessary at specific times in the execution of the program in order for the program to execute in an efficient manner; in contrast, since writes to memory are merely writing the result of the program execution, the actual time at which the writing occurs is not as critical since the execution of later instructions may not depend upon the result.\nBy way of further background, write buffers have been provided in microprocessors, such write buffers being logically located between on-chip cache memory and the bus to main memory. These conventional post-cache write buffers receive data from the cache for a write-through or write-back operation; the contents of the post-cache write buffer are written to main memory under the control of the bus controller, at times when the bus becomes available.\nBy way of further background, some pipelined architecture microprocessors operate according to speculative execution in order to maintain the pipeline full despite conditional branch or jump instructions being present in the program sequence. Speculative execution requires that predictive branching be performed, where the microprocessor predicts whether the conditional branch will be taken or not taken according to an algorithm; the predicted path is then executed in the pipeline. It is important that the results of speculative executed instructions not be written to memory or cache, because if the prediction is incorrect, it may be difficult or impossible to recover from the incorrectly performed memory write.\nAnother type of situation can occur where instructions are processed in a pipeline, including writes to memory, where an earlier instruction has an exception condition (e.g., divide-by-zero) for which the program execution should be immediately stopped.\nIt is an object of the present invention to provide a microprocessor architecture which buffers the writing of data from the CPU core into a write buffer, prior to retiring of the data to a cache, and in which recovery from speculative execution or exceptions can be readily performed.\nIt is a further object of the present invention to provide such an architecture which prevents the writing of data to memory during a speculative execution sequence.\nIt is a further object of the present invention to provide such an architecture which allows for multiple degrees of speculative execution.\nOther objects and advantages of the present invention will be apparent to those of ordinary skill in the art having reference to the following specification in combination with the drawings."}
{"text": "In recent years, a sputter neutral particle mass spectrometry apparatus using a focused ion beam device and a light beam oscillation device has been developed. In light beam post-ionized neutral particle mass spectrometry using a focused ion beam, an ion beam is irradiated on a sample to generate neutral particles, and a light beam is made incident horizontally on the surface of the sample to perform post-ionization a single time. Here, a secondary ion mass spectrometry apparatus is equipped with a light beam for post-ionization to improve detection sensitivity. Therefore, the measuring system is simplified to facilitate optimization of the timing for irradiating the light beam and the timing for drawing in ions.\nIn such neutral particle mass spectrometry, since the output of the light beam is small, it has been necessary to make the light beams converge at a particular position on the sample surface in order to secure sufficient detection sensitivity. Therefore, to prevent background noise caused by direct contact of the light beams with the sample from occurring, the optical axis of the light beam has been required to be designed horizontal to, and at a fixed distance from the sample surface. As a result, the light beam irradiation timing needed to be delayed, and thus density per unit volume of a sputter neutral particle group decreased, causing a decline in yield.\nIn recent years, objects to be analyzed by the sputter neutral particle mass spectrometry apparatus have become microscopic. Therefore, there is a need for a sputter neutral particle mass spectrometry apparatus which is capable of improving detection sensitivity by increasing the yield of any element in the region of analysis."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a network of devices. More particularly, it relates to an apparatus and system for coupling and decoupling data storage initiator devices to a network without disrupting the network.\n2. Related Art\nThe data storage market includes a number of vendors and products. Unfortunately, integrating various products from different vendors is difficult, and it requires a substantial investment due to a lack of interoperability standards.\nIn one instance, in order to increase system performance and lower costs, the manufacturers of blade servers and other storage devices are considering integrating a Fibre Channel fabric switch into their devices. However, a blade server with an integrated fabric switch is likely to have difficulties communicating to an external network because of incompatibilities and proprietary features. The conventional wisdom is that such devices are connected using a Fibre Channel E-Port or B-Port topology, thereby allowing fabric related information to be communicated. But, this causes many currently available Fibre Channel fabric switches to be reconfigured to a mode in which proprietary features are turned off and functions are disabled, resulting in a disruption of the network. It is also likely to create network management problems.\nWhat is needed, therefore, are new ways for integrating products from different venders to an existing network that overcome the deficiencies noted above."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for controlling multiple powers for sequential supply and reversal interruption and, more particularly, to an apparatus for controlling multiple powers which is capable of sequentially turning on or reversely turning off the multiple powers in their priorities for systems or components to be supplied with the multiple powers such as a liquid crystal display (LCD) module.\n2. Discussion of Related Art\nHereinafter, the supply and the interruption have the same meaning as ON and OFF, which will be used together for simplicity of description.\nIn general, multiple powers should be supplied to components or systems requiring the multiple powers for safe drive in a sequential manner to comply with their priorities of supplying the powers based on the respective specifications, and should be interrupted in its reversal order. To detail this, in the case of the power supply for driving the LCD module, four kinds of power supply voltage, 3.3V, 5V, −9V, and 15V, are sequentially supplied to the LCD module by respective time intervals of 10 ms, image data are then applied thereto. In addition, in the case of interrupting the power supply, it is interrupted in its reversal order.\nFor example, one of methods capable of performing sequential supply and reversal interruption on the multiple powers in accordance with the conventional art will be described. FIG. 1 is a circuit view of an apparatus for controlling multiple powers in accordance with the conventional art.\nThe apparatus for controlling the multiple powers according to the conventional art as shown in FIG. 1 is configured to sequentially switch the multiple power supply voltages by sequentially driving relays using a silicon controlled rectifier (SCR) and counter logical elements. In this configuration, a pair of counter elements, i.e., a counter element for sequential turn on and a counter element for reversal turn off are required to generate a control signal for switching one power supply voltage. An upper block including the counter element is responsible for sequential on control and a lower block is responsible for sequential off control. In addition, one SCR and two relays are required in its output terminal in order to turn off one voltage having levels different from each other in the above-described configuration.\nAccordingly, in the apparatus for controlling the multiple powers according to the conventional art, the more the voltages to be controlled are present, the more logical gates and counter semiconductor elements constituted by flip-flops or latches should be used and many relays should also be used, which cause volume to be increased and control to be complicated."}
{"text": "A recently developed type of secondary or rechargeable electrical conversion device comprises: (1) one or more anodic reaction zones containing a molten alkali metal anode-reactant, e.g., sodium, in electrical contact with an external circuit; (2) one or more cathodic reaction zones containing (a) a cathodic reactant comprising a liquid electrolyte, e.g., sulfur or a mixture of sulfur and molten polysulfide, which is electrochemically reversibly reactive with said anodic reactant, and (b) a conductive electrode which is at least partially immersed in said cathodic reactant; and (3) a solid electrolyte comprising a cationpermeable barrier to mass liquid transfer interposed between and in contact with said anodic and cathodic reaction zones. As used herein the term \"reactant\" is intended to mean both reactants and reaction products.\nDuring the discharge cycle of such a device, molten alkali metal atoms such as sodium surrender an electron to an external circuit and the resulting cation passes through the solid electrolyte barrier and into the liquid electrolyte to unite with polysuflfide ions. The polysulfide ions are formed by charge transfer by reaction of the cathodic reactant with electrons conducted through the electrode from the external circuit. Because the ionic conductivity of the liquid electrolyte is less than the electronic conductivity of the electrode material, it is desirable during discharge that both electrons and sulfur be applied to and distributed along the surface of the conductive electrode in the vicinity of the cation-permeable solid electrolyte. When sulfur and electrons are so supplied, polysulfide ions can be formed near the solid electrolyte and the alkali metal cations can pass out of the solid electrolyte into the liquid electrolyte and combine to form alkali metal polysulfide near the solid electrolyte.\nDuring the charge cycle of such a device when a negative potential larger than the open circuit cell voltage is applied to the anode the opposite process occurs. Thus electrons are removed from the alkali metal polysulfide by charge transfer at the surface of the electrode and are conducted through the electrode to the external circuit, and the alkali metal cation is conducted through the liquid electrolyte and solid electrolyte to the anode where it accepts an electron from the external circuit. Because of the aforementioned relative conductivities of the ionic and electronic phases, this charging process occurs preferentially in the vicinity of the solid electrolyte and leaves behind molten elemental sulfur. As can be readily appreciated the production of large amounts of sulfur near the surface of the cation-permeable membrane has a limiting effect on rechargeability. This is the case since sulfur is nonconductive and when it covers surfaces of the electrode, charge transfer is inhibited and the charging process is greatly hindered or terminated. Thus, in order to improve the rechargeability of a cell of this type it is necessary not only to supply polysulfide to the surface of the electrode in the vicinity of the cation-permeable membrane, but also to remove sulfur therefrom.\nU.S. Pat. No. 3,811,943 and United States patent application Ser. No. 545,048 filed Jan. 29, 1975 now U.S. Pat. No. 3,980,496 both disclose energy conversion device designs which allow or promote improved mass transportation of reactants and reaction products to and from the vicinity of the solid electrolyte and the porous electrode during both discharge and charge. In the device disclosed in the patent an ionically conductive solid electrolyte is located between a first reactant in one container and a second reactant in another container. An electrode for one of the reactants comprises a layer of porous, electronically conductive material having one surface in contact with one side of the ionically conductive solid electrolyte and the other surface in contact with a structurally integral electronically conductive member permeable to mass flow of its reactants and electrically connected to the external circuit. An open volume exists between the structurally integral conductive member and the container wall to promote free flow and mixing of the reactants. Reactants also flow readily through the conductive member into the layer of porous electronically conductive material. The conductive member distributes electrons to the porous, conductive material which in turn transfers electrons to or from the reactants.\nThe improvement disclosed in Ser. No. 545,048 comprises designing the cathodic reaction zone of the device such that there are a plurality of channels and/or spaces within said zone which are free of porous conductive electrodes and which are thus adapted to allow free flow of the molten cathodic reactants during operation of the device. This flow results from free convection within the channels and/or spaces, and from wicking of cathodic reactants within the conductive porous material.\nThe secondary battery or cell designs disclosed and claimed in the aforementioned U.S. patent and Ser. No. 545,048 are effective in promoting distribution of reactants during both charge and discharge, thereby increasing the capacity of the device. However, the capacity of the device is limited by the amount of anodic and cathodic reactant available for reaction.\nIt is the object of this invention to provide an improved secondary battery or cell which is effective in promoting distribution of reactants during charge and discharge and which provides for an increased supply of reactants, thereby increasing the ampere-hour capacity of the device."}
{"text": "There currently exists a variety of absorbent articles that have been developed to improve the health and hygiene of the user. For example, infant diapers, training pants, adult incontinence articles, feminine care articles, and the like are commonly used to absorb urine, menses, and other bodily exudates. However, despite the useful nature of these products, some users desire greater discretion regarding their need for particular absorbent articles. This may be particularly relevant to some adult users with incontinence. As such, some users may obtain absorbent articles by home delivery to avoid the embarrassment or self-consciousness of shopping in a traditional grocery store or drug store.\nHowever, even with home delivery, some users are concerned about the cases of branded product being seen by others during delivery, storage, and/or disposal. To address this need, some manufacturers of absorbent articles ship the products in cases having no visible branding. This addresses the need for discretion but can cause issues with some retailers that want to see the branding on the cases to ensure proper merchandise inventory, transfer, and/or stocking. In order to serve both needs, some manufactures have provided two different case designs to satisfy the needs of the two different customers (i.e., branded cases for the retailers and unbranded cases for home delivery). However, this solution creates logistical and economic inefficiencies by requiring the manufacturer of the absorbent articles to manage two different case designs for a single product code.\nOther methods of addressing discretion have included manufactures of absorbent articles providing the same branded cases to both retailers and home delivery providers. In turn, some of the home delivery providers have emptied the cases, turned the cases inside out, and refilled the cases such that the brand is printed on the inside of the case. Transforming the cases in this way eliminates the need for two different case designs but requires additional labor and expense on the part of the home delivery providers. Furthermore, this solution does not address the need for discretion when storing or disposing of the branded case because the branding is still visible on the inner surface of the case.\nIn other situations when manufactures of absorbent articles provide the same branded cases to both retailers and home delivery providers, the home delivery providers remove the product from the branded cases and repack the product into unbranded cases. This solution addresses the need for discretion during receipt, storage, and disposal of the case but re-introduces the problem of having two different case designs (i.e., branded and unbranded). The only difference is that the expense has transferred from the manufacturer of the absorbent articles to the home delivery provider.\nThus, there exists a need for a single case design that meets the branding requirements of retailers and the discretion requirements of home delivery providers. There also exists a need for a method to use the single case design to satisfy the needs of all customers/retailers."}
{"text": "The present invention relates generally to communication networks, and more specifically to the monitoring and management of quality of signal transmissions in a communications network.\nOriginally, cable networks were established to transmit television signals to homes and offices. Cable networks provided advantages over transmission television networks that included providing a clearer signal and a greater selection of channels. These networks were made up of co-axial cables routed in a tree and branch structure to customer sites and were intended simply to provide customers with analog television signals.\nMore recently, cable networks have been converted to transmit digital signals in a hybrid of fiber optic cable and co-axial cable structures. These converted networks accommodate not only traditional, analog television signals but also digital television signals, digital data signals and telephone signals. These cable networks use packet switching techniques to transfer information (data) over packet networks using packet switching techniques.\nDigital television signals provide a crisper, more detailed picture along with enhanced sound. With the capability to transmit digital data signals, cable networks may now be coupled to the Internet thereby providing homes and offices access to the Internet. This Internet access is generally faster than access provided by other technologies. In addition, cable networks may now be used to transmit telephone voice signals in the form of packetized data. This involves a substantial savings in the amount of wiring around houses and offices because one co-axial cable can now carry analog, digital and voice signals.\nAnother signal transmission technology that offers significant improvement in data transfer is the Digital Subscriber Line (DSL) technology. DSL technology provides increased communications bandwidth while using existing twisted-pair copper lines that are prevalent throughout much of the world. DSL delivers a basic data transfer rate of 128 kbps. High speed DSL, or HDSL, can deliver a data transfer rate of 1.544 megabits per second (Mbps) in North America, and 2.048 Mbps elsewhere. Asymmetric DSL, or ADSL, can deliver data rates ranging from 1.5 to 9.0 Mbps on a downstream or receiving path, and 16 to 800 kbps on an upstream or sending path. Taken together, varying DSL technologies are referred to as xDSL.\nA conventional xDSL communication network includes a Main Distribution Frame (MDF), an access matrix, a DSL Access Multiplexer (DSLAM), and a test unit. The MDF is coupled to the access matrix, which itself is coupled to the DSLAM and the test unit. The MDF, the access matrix, the test unit, and the DSLAM each reside at an xDSL service provider site (or Central Office). At a customer site, a set of Customer Premises Equipment (CPE) units is connected to the MDF. Each CPE unit includes an xDSL modem.\nTo monitor the quality of the transmission of telephone or voice signals in the form of packetized data, various transmission impairment tests (TIT) are used. One common type of transmission impairment test is the Perceptual Speech Quality Measurement (PSQM) according to the International Telecommunication Union (ITU) P.861 Standard. Information about ITU and ITU Standard P.861 is available at http://www.itu.ch. Determination of a PSQM score typically involves sending of a PSQM file from a home device to a remote device, and calculating a PSQM score by the remote device based on the PSQM file received at the remote device. A PSQM file is a file including digital signals representative of voice signals of various people (e.g., men, women, children, etc.) speaking different languages. Alternatively, a PSQM score can be determined by sending a PSQM file from a remote device to a home device, and calculating the PSQM score by the home device based on the PSQM file received at the home device. The PSQM score is a number between one and seven where a higher number represents higher quality.\nAnother type of TIT is the Perceptual Evaluation of Speech Quality (PESQ) according to the ITU P.862 Standard. Information about ITU P.862 is available at http://www.itu.ch. PESQ has been found to be more accurate than PSQM at predicting quality in a very wide range of networks, including the speed transmission quality of packet-oriented networks. The PESQ score is a number between one and seven where a higher number represents higher quality.\nThe quality of the transmission of telephone or voice signals in the form of packetized data is also measured by another criteria commonly known as the Quality of Services (QoS). The QoS is a measurement of the amount of packet losses, jitter, delay, etc in the signal transmission. Transmission impairment tests (e.g., PSQM and PESQ), and QoS are important criteria to determine the quality of signal transmissions in a packetized network.\nIt should be noted that the TIT scores is more difficult to obtain than the QoS score as the TIT score, such as the PSQM score, varies with time and location of the measurement. Therefore, it would be desirable to correlate the TIT scores with a corresponding QoS score by measuring the TIT scores and the QoS score for the same signal transmission. This would allow a user (or an administrator) of the packet network to use a QoS score to predict the corresponding TIT scores. When the predicted TIT score drops below a minimum value indicative of network problems, the administrator will be informed to determine whether services to the network are needed to restore the signal transmission quality.\nHowever, there are problems associated with the prior art systems that make the correlation of TIT scores and QoS scores complicated and expensive. First, even for a small region, thousands of phone calls are being made at any given time which makes it very difficult to keep track of a particular test telephone call to allow the calculation of TIT score and QoS score. Second, the TIT scores vary from location to location and also from time to time for the same location because of variations in call patterns and the traffic conditions of the packet network.\nThus, there has long been a need for a more economical and simple monitoring method and system, which would monitor and calculate TIT scores, such as PSQM scores and PESQ scores, and QoS scores for signal transmissions over a packet network; and then use QoS scores to provide an accurate prediction of the corresponding TIT scores.\nThe present invention provides a method and system for monitoring the quality of signal transmissions in a communications network.\nThe present invention further provides an economical and simple method and system for calculating both Transmission Impairment Test (TIT) scores and QoS scores for signal transmissions over a communications network.\nThe present invention still further provides a method and system for providing a correlation between TIT scores and QoS scores of signal transmissions over a communications network, thus allowing a user of the communications network to use a QoS score to predict the corresponding TIT score. When the TIT score falls below a minimum TIT score, the user may be informed and will have an opportunity to determine if service to the communications network is required to restore signal transmission quality.\nThe present invention provides a communications network having a cable modem termination system (CMTS); a voice band tester (VBT) coupled to the CMTS; a cable modem tester coupled to the CMTS; and a Voice over Internet Packet (VoIP) monitoring device coupled to the CMTS and the VBT. The VBT is located at a first location and the cable modem tester is located at a second location remote from the first location. The cable modem tester is adapted to provide a first communication signal to the voice band tester via the CMTS. The VoIP monitoring device is adapted to monitor the first communication signal and calculate a first Quality of Services (QoS) score based on traffic density between the CMTS and the VBT. The VBT is further adapted to calculate a first Transmission Impairment Test (TIT) score based on the first communication signal and a first received communication signal received by the VBT from the cable modem tester and provide the first TIT score to the VoIP monitoring device. The first TIT score could be a score such as a Perceptual Speech Quality Measurement (PSQM) score, or a Perceptual Evaluation of Speech Quality (PESQ) score.\nThe present invention further provides a method for monitoring quality of signal transmissions within a communications network. The method includes the steps of: (a) providing a first communication signal from a cable modem tester located at a first location to a voice band tester (VBT) located at a second location remote from the first location via a Cable Modem Termination System (CMTS); (b) identifying the first communication signal and begins monitoring signal transmissions from the cable modem tester to the VBT via the CMTS; (c) calculating a first Transmission Impairment Test (TIT) score based on the first communication signal and a first received communication signal received by the VBT from the cable modem tester; (d) providing the first TIT score to a Voice over Internet Packet (VoIP) monitoring device; and (e) calculating a first Quality of Services (QoS) score based on traffic density between the CMTS and the VBT.\nThe above and additional advantages of the present invention will become apparent to those skilled in the art from a reading of the following detailed description when taken in conjunction with the accompanying drawings."}
{"text": "In thermal dye transfer printing, an image is formed on a receptor sheet by selectively transferring an image forming material to the receptor sheet from a dye donor sheet. Material to be transferred from the dye donor sheet is selected by a thermal printhead, which consists of small, electrically heated elements. These elements transfer image-forming material from the dye donor sheet to areas of the dye receptor sheet in an imagewise manner.\nThere are three broad classes of thermal transfer systems that are known, (1) chemical reaction systems, (2) thermal mass transfer systems, and (3) thermal dye transfer systems.\nIn chemical reaction systems, the image is formed upon the receptor as a result of the imagewise transfer of some chemical reactant from the donor sheet. An example is the transfer of a mobile molecule, such as phenol, to the receptor sheet, which bears a leuco compound thereon. The phenol is transferred by being volatilized by the heat from the thermal print head, and, upon reaching the receptor sheet, reacts with the leuco compound to convert it from the colorless to the colored form. Alternately, the phenol can be on the receptor sheet and the leuco compound can be on the donor sheet.\nIn thermal mass systems, no color forming chemical reaction takes place. Instead, the image is formed simply by the transfer of a mass of material containing colorant therein, such as pigment-filled polymer coatings.\nIn thermal dye transfer systems, a dye donor sheet is used in combination with a dye receptor sheet wherein, with the application of heat, a dye is transferred onto the receptor sheet at a controlled amount to obtain a dye image having gradation like in a photograph.\nEach system has its own advantages and disadvantages for the particular application of thermal printing. Various problems have been encountered with each proposed system. For thermal dye transfer, dye release layers have been proposed to enable efficient transfer of the dye layer from the dye donor sheet. Also, various dye-permeable release layers on the dye receptor layer have been proposed. The dye-permeable release layer is coated over the dye receptor layer, and is formulated to prevent sticking between the donor layer and the receptor layer during the transfer of the dye across the binder membranes. The release layer must also be formulated to allow effective transfer of the dye through the release layer. In general, many of these problems have been related to a specific resin used in the composition of the dye donor or dye receptor layers.\nSelection of the functional resin systems for the dye donor and the receptor sheet layers has been the topic of concern for many proposed dye systems. In consideration of the above mentioned requirements, efficient dye transfer and sticking between the dye donor layer and the dye receptor layer during transfer, a good, functional resin system to eliminate some of these problems is needed. Interestingly, a unique resin system has been found that provides an efficient working dye donor sheet and dye receptor sheet. The resin system comprises a chlorinated polyvinyl chloride (CPVC).\nU.S. Pat. No. 3,584,576 describes a heat sensitive stencil sheet comprising a film adhered to a porous thin fibrous sheet. The stencil sheet is perforated by exposure to infrared rays. The film consists essentially of at least 75% by weight of a chlorinated polyvinyl chloride resin, the balance being a polyvinyl chloride resin. A colorant may be present in the film. Upon being heated by infrared radiation, the film melts and forms perforations. The pores in the remaining fibrous sheet enable stencilling to be done through the perforations and the sheet.\nThere are noticable difference between the above mentioned prior art use of CPVC in a thermally sensitive stencil applications and in the present invention. The prior art uses CPVC merely as a resinous binder, with or without other resinous binders. It is, in particular, used not as a receptor layer for a thermal dye transfer sheet, but as the thermoplastic binder for a thermal stencil sheet. The novel use of the CPVC resin in the thermal dye transfer printing of the present invention has been found to give surprisingly new and unique properties for use as the primary resinous thermal plastic binder in both a dye donor sheet, and a dye receptor sheet. Typically, commercially available dye donor sheets and dye receptor sheets are comprised of chemically different binders with different functionals."}
{"text": "Semicarbazide-sensitive amine oxidase (SSAO) activity is an enzyme activity expressed by Vascular Adhesion Protein-1 (VAP-1) or Amine Oxidase, Copper Containing 3 (AOC3), belongs to the copper-containing amine oxidase family of enzymes (EC.1.4.3.6). Therefore inhibitors of the SSAO enzyme may also modulate the biological functions of the VAP-1 protein. Members of this enzyme family are sensitive to inhibition by semicarbazide and utilize cupric ion and protein-derived topa quinone (TPQ) cofactor in the oxidative deamination of primary amines to aldehydes, hydrogen peroxide, and ammonia according to the following reaction:R—CH2—NH2+O2→R—CHO+H2O2+NH3 \nKnown substrates for human SSAO include endogenous methylamine and aminoacetone as well as some xenobiotic amines such as benzylamine [Lyles, Int. J. Biochem. Cell Biol. 1996, 28, 259-274; Klinman, Biochim. Biophys. Acta 2003, 1647(1-2), 131-137; Matyus et al., Curr. Med. Chem. 2004, 11(10), 1285-1298; O'Sullivan et al., Neurotoxicology 2004, 25(1-2), 303-315]. In analogy with other copper-containing amine oxidases, DNA-sequence analysis and structure determination suggest that the tissue-bound human SSAO is a homodimeric glycoprotein consisting of two 90-100 kDa subunits anchored to the plasma membrane by a single N-terminal membrane spanning domain [Morris et al., J. Biol. Chem. 1997, 272, 9388-9392; Smith et al., J. Exp. Med. 1998, 188, 17-27; Airenne et al., Protein Science 2005, 14, 1964-1974; Jakobsson et al., Acta Crystallogr. D Biol. Crystallogr. 2005, 61(Pt 11), 1550-1562]. SSAO activity has been found in a variety of tissues including vascular and non-vascular smooth muscle tissue, endothelium, and adipose tissue [Lewinsohn, Braz. J. Med. Biol. Res. 1984, 17, 223-256; Nakos & Gossrau, Folia Histochem. Cytobiol. 1994, 32, 3-10; Yu et al., Biochem. Pharmacol. 1994, 47, 1055-1059; Castillo et al., Neurochem. Int. 1998, 33, 415-423; Lyles & Pino, J. Neural. Transm. Suppl. 1998, 52, 239-250; Jaakkola et al., Am. J. Pathol. 1999, 155, 1953-1965; Morin et al., J. Pharmacol. Exp. Ther. 2001, 297, 563-572; Salmi & Jalkanen, Trends Immunol. 2001, 22, 211-216]. In addition, SSAO protein is found in blood plasma and this soluble form appears to have similar properties as the tissue-bound form [Yu et al., Biochem. Pharmacol. 1994, 47, 1055-1059; Kurkijarvi et al., J. Immunol. 1998, 161, 1549-1557]. It has recently been shown that circulating human and rodent SSAO originates from the tissue-bound form [Gökürk et al., Am. J. Pathol. 2003, 163(5), 1921-1928; Abella et al., Diabetologia 2004, 47(3), 429-438; Stolen et al., Circ. Res. 2004, 95(1), 50-57], whereas in other mammals the plasma/serum SSAO is also encoded by a separate gene called AOC4 [Schwelberger, J. Neural. Transm. 2007, 114(6), 757-762].\nThe precise physiological role of this abundant enzyme has yet to be fully determined, but it appears that SSAO and its reaction products may have several functions in cell signalling and regulation. For example, recent findings suggest that SSAO plays a role in both GLUT4-mediated glucose uptake [Enrique-Tarancon et al., J. Biol. Chem. 1998, 273, 8025-8032; Morin et al., J. Pharmacol. Exp. Ther. 2001, 297, 563-572] and adipocyte differentiation [Fontana et al., Biochem. J. 2001, 356, 769-777; Mercier et al., Biochem. J. 2001, 358, 335-342]. In addition, SSAO has been shown to be involved in inflammatory processes where it acts as an adhesion protein for leukocytes [Salmi & Jalkanen, Trends Immunol. 2001, 22, 211-216; Salmi & Jalkanen, in “Adhesion Molecules: Functions and Inhibition” K. Ley (Ed.), 2007, pp. 237-251], and might also play a role in connective tissue matrix development and maintenance [Langford et al., Cardiovasc. Toxicol. 2002, 2(2), 141-150; Goktark et al., Am. J. Pathol. 2003, 163(5), 1921-1928]. Moreover, a link between SSAO and angiogenesis has recently been discovered [Noda et al., FASEB J. 2008, 22(8), 2928-2935], and based on this link it is expected that inhibitors of SSAO have an anti-angiogenic effect.\nSeveral studies in humans have demonstrated that SSAO activity in blood plasma is elevated in conditions such as congestive heart failure, diabetes mellitus, Alzheimer's disease, and inflammation [Lewinsohn, Braz. J. Med. Biol. Res. 1984, 17, 223-256; Boomsma et al., Cardiovasc. Res. 1997, 33, 387-391; Ekblom, Pharmacol. Res. 1998, 37, 87-92; Kurkijarvi et al., J. Immunol. 1998, 161, 1549-1557; Boomsma et al., Diabetologia 1999, 42, 233-237; Meszaros et al., Eur. J. Drug Metab. Pharmacokinet. 1999, 24, 299-302; Yu et al., Biochim. Biophys. Acta 2003, 1647(1-2), 193-199; Mntyus et al., Curr. Med. Chem. 2004, 11(10), 1285-1298; O'Sullivan et al., Neurotoxicology 2004, 25(1-2), 303-315; del Mar Hernandez et al., Neurosci. Lett. 2005, 384(1-2), 183-187]. The mechanisms underlying these alterations of enzyme activity are not clear. It has been suggested that reactive aldehydes and hydrogen peroxide produced by endogenous amine oxidases contribute to the progression of cardiovascular diseases, diabetic complications and Alzheimer's disease [Callingham et al., Prog. Brain Res. 1995, 106, 305-321; Ekblom, Pharmacol. Res. 1998, 37, 87-92; Yu et al., Biochim. Biophys. Acta 2003, 1647(1-2), 193-199; Jiang et al., Neuropathol Appl Neurobiol. 2008, 34(2), 194-204]. Furthermore, the enzymatic activity of SSAO is involved in the leukocyte extravasation process at sites of inflammation where SSAO has been shown to be strongly expressed on the vascular endothelium [Salmi et al., Immunity 2001, 14(3), 265-276; Salmi & Jalkanen, in “Adhesion Molecules: Functions and Inhibition” K. Ley (Ed.), 2007, pp. 237-251]. Accordingly, inhibition of SSAO has been suggested to have a therapeutic value in the prevention of diabetic complications and in inflammatory diseases [Ekblom, Pharmacol. Res. 1998, 37, 87-92; Salmi et al., Immunity 2001, 14(3), 265-276; Salter-Cid et al., J. Pharmacol. Exp. Ther. 2005, 315(2), 553-562].\nSSAO knockout animals are phenotypically overtly normal but exhibit a marked decrease in the inflammatory responses evoked in response to various inflammatory stimuli [Stolen et al., Immunity 2005, 22(1), 105-115]. In addition, antagonism of its function in wild type animals in multiple animal models of human disease (e.g. carrageenan-induced paw inflammation, oxazolone-induced colitis, lipopolysaccharide-induced lung inflammation, collagen-induced arthritis, endotoxin-induced uveitis) by the use of antibodies and/or small molecules has been shown to be protective in decreasing the leukocyte infiltration, reducing the severity of the disease phenotype and reducing levels of inflammatory cytokines and chemokines [Kirton et al., Eur. J. Immunol. 2005, 35(11), 3119-3130; Salter-Cid et al., J. Pharmacol. Exp. Ther. 2005, 315(2), 553-562; McDonald et al., Annual Reports in Medicinal Chemistry 2007, 42, 229-243; Salmi & Jalkanen, in “Adhesion Molecules: Functions and Inhibition” K. Ley (Ed.), 2007, pp. 237-251; Noda et al., FASEB J. 2008 22(4), 1094-1103; Noda et al., FASEB J. 2008, 22(8), 2928-2935]. This anti-inflammatory protection seems to be afforded across a wide range of inflammatory models all with independent causative mechanisms, rather than being restricted to one particular disease or disease model. This would suggest that SSAO may be a key nodal point for the regulation of the inflammatory response, and it is therefore likely that SSAO inhibitors will be effective anti-inflammatory drugs in a wide range of human diseases. SSAO (VAP-1) is up regulated in gastric cancer and has been identified in the tumour vasculature of human melanoma, hepatoma and head and neck tumours (Yoong K F, McNab G, Hubscher S G, Adams D H. (1998), J Immunol 160, 3978-88; Irjala H, Salmi M, Alanen K, Gre'nman R, Jalkanen S (2001), Immunol. 166, 6937-6943; Forster-Horvath C, Dome B, Paku S, et al. (2004), Melanoma Res. 14, 135-40.). One report (Marttila-Ichihara F, Castermans K, Auvinen K, Oude Egbrink M G, Jalkanen S, Griffioen A W, Salmi M. (2010), J Immunol. 184, 3164-3173.) has shown that mice bearing enzymically inactive VAP-1 grow melanomas more slowly, and have reduced tumour blood vessel number and diameter. The reduced growth of ic) these tumours was also reflected in the reduced (by 60-70%) infiltration of myeloid suppressor cells. Encouragingly VAP-1 deficiency had no effect on vessel or lymph formation in normal tissue.\nSmall molecules of different structural classes have previously been disclosed as SSAO inhibitors, for example in WO 02/38153 (tetrahydroimidazo[4,5-c]pyridine derivatives), in WO 03/006003 (2-indanylhydrazine derivatives), in WO 2005/014530 (allylhydrazine and hydroxylamine (aminooxy) compounds) and in WO 2007/120528 (allylamino compounds). Additional SSAO inhibitors are disclosed in PCT/EP2009/062011 and PCT/EP2009/062018.\nOur co-pending International Patent Application No.: PCT/EP2011/053818 relates to SSAO inhibitors of formula (I) or a pharmaceutically acceptable salt, or N-oxide thereof:\nwhereinR1 is phenyl or 6-membered heteroaryl, optionally substituted with one or more substituents selected from halogen, cyano, C1-4-alkyl, halo-C1-4-alkyl, C1-4alkoxy-C1-4alkyl, hydroxy-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)amino-C1-4-alkyl, —NR4AR4B, —NR6C(O)OR5, —NR6C(O)R5, —NR6C(O)NR4AR4B, —C(O)NR4AR4B, —C(O)R5, —C(O)OR5, and —NR6S(O)2R5;A is a bond;R2 is —B-Q-[R3]n or —B—R3;wherein n=1, 2, 3, or 4B is a bond, O, NR4, —C(O)— or C1-3-alkylene;Q is saturated or partially unsaturated monocyclic 3-7 membered heterocyclic or C3-7-cycloalkyl ring;when R2 is —B-Q-[R3]n, R3 is selected from hydrogen, halogen, cyano, amino, hydroxyl, oxo, C1-4-alkyl, halo-C1-4-alkyl, C1-4-alkoxy, C1-4alkoxy-C1-4alkyl, hydroxy-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)amino-C1-4-alkyl, —NR4AR4B, —NR6C(O)OR5, —NR6C(O)R5, —NR6C(O)NR4AR4B, —C(O)NR4AR4B, —C(O)R5, —C(O)OR5, —NR6S(O)2R5, —S(O)2R5, phenyl-C1-4-alkyl and heteroaryl-C1-4-alkyl, and wherein any phenyl or heteroaryl residue is optionally substituted with one or more substituents selected from halogen, cyano, C1-4-alkyl, halo-C1-4-alkyl, C1-4-alkoxy, C1-4alkoxy-C1-4alkyl, hydroxy-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)amino-C1-4-alkyl, —NR4AR4B, —NR6C(O)OR5, —NR6C(O)R5, —NR6C(O)NR4AR4B, —C(O)NR4AR4B, —C(O)R5, —C(O)OR5, —NR6S(O)2R5;when R2 is —B—R3, R3 is selected from amino, C1-4-alkoxy, C1-4alkoxy-C1-4alkyl, hydroxy-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)amino-C1-4-alkyl, —NR4AR4B, —NR6C(O)OR5, —NR6C(O)R5, —NR6C(O)NR4AR4B, —C(O)NR4AR4B, —C(O)R5, —NR6S(O)2R5, phenyl-C1-4-alkyl and heteroaryl-C1-4-alkyl, and wherein any phenyl or heteroaryl residue is optionally substituted with one or more substituents selected from halogen, cyano, C1-4-alkyl, halo-C1-4-alkyl, C1-4-alkoxy, C1-4alkoxy-C1-4alkyl, hydroxy-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)amino-C1-4-alkyl, —NR4AR4B, —NR6C(O)OR5, —NR6C(O)R5, —NR6C(O)NR4AR4B, —C(O)NR4AR4B, —C(O)R5, —C(O)OR5, —NR6S(O)2R5, provided that when R2 is —B—R3, and B is a bond and R3 is —C(O)R5, then R5 is not hydrogen;R4A, R4B and R5 are each independently selected from hydrogen, C1-4-alkyl, hydroxy-C1-4-alkyl, halo-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)-amino-C1-4-alkyl or C1-4alkoxy-C1-4alkyl; or R4A and R4B together with the nitrogen to which they are attached form a cyclic amino group such as a piperidinyl, piperazinyl, N-substituted piperazinyl, morpholinyl or homopiperidinyl group;R4A1 is selected from C1-4-alkyl, hydroxy-C1-4-alkyl, halo-C1-4-alkyl, cyano-C1-4-alkyl, amino-C1-4-alkyl, C1-4-alkylamino-C1-4-alkyl, di(C1-4-alkyl)amino-C1-4-alkyl or C1-4alkoxy-C1-4alkyl; or R4A1 and R4B together with the nitrogen to which they are attached form a cyclic amino group such as a piperidinyl, piperazinyl, N-substituted piperazinyl, morpholinyl or homopiperidinyl group;R6 is hydrogen or C1-4-alkyl; andX is selected from the radicals of formulae (1-16) wherein the bond marked * is attached to R1A- and the bond marked ** is attached to —R2:\nwherein Y is selected from hydrogen, hydroxyl, amino, —NHR6, —OCH3;Z is selected from hydrogen, fluorine, hydroxyl, C1-4-alkoxy, halo-C1-4-alkyl, CONH2, cyano, SO2NH2, amino, —NHR6;W is selected from H, C1-4-alkyl, halo-C1-4-alkyl;"}
{"text": "One of the major hypotheses that has recently been proposed for the cause of the observed effect of particulate matter on human health is that high numbers of ultrafine particles (for example, particles, diameters less than 0.1 μm) are more problematic than the particle mass that is now the basis of the National Ambient Air Quality Standards. For example, significant associations of elevated cardiovascular and respiratory disease mortality with various fine (and ultrafine) particle indices have been found in one study based in Erfurt, Germany. Specifically, significant associations were found between mortality and ultrafine particle number (NC), ultrafine particle mass (MC), fine particle mass or SO2 concentrations. The correlation between MC 0.01-2.5 and NC 0.01-0.1 is only moderate, suggesting that it may be possible to partially separate effects of ultrafine and fine particles. Thus, measurements of the ultrafine particle concentrations as well as particle mass are needed to help provide more data to examine these relationships.\nIn addition, there are a variety of particle counting needs in industrial settings. With the increased emphasis on nanometer sized particles for the production of nanostructured materials, particle counters can be important in process control. There are currently only a limited number of instruments available to make such measurements.\nOne method or device that has shown to be an effective means of detecting such fine particles is the use of heterogeneous nucleation with a turbulent mixing condensation nuclei counter (that is, a “TMCNC”). In a TMCNC system, a gas (for example, air) containing the particles to be measured is turbulently mixed with a stream of air saturated with a condensable vapor so that the vapor cools and condensates, that is nucleates, onto the particles. The resulting droplets or nuclei then grow to a size whereby they can be effectively detected by, for example, light scattering. Although the concept of a TMCNC has been available for over 15 years (see for example, McMurry, 2000), the TMCNC has not been developed into a viable commercial instrument capable of particle detection down to 2 nanometers.\nThe Condensation Nuclei Counter (CNC), which grows primary particles (nuclei) up to a more easily detectable size, is one of the most widely used devices for studying particles below 0.1 micrometer (μm). A general description of CNC's is given in many books and reviews (for example, see Willeke and Baron, 1993 and McMurry, 2000). Several types of CNCs are used in researching aerosols, that is, suspensions of fine solid or liquid particles in a gas, typically, air. The main difference among these CNC designs is the way the devices produce supersaturation that leads to particle growth up to a predetermined size for subsequent detection. In an expansion-type CNC, supersaturation is generated by adiabatic cooling during pressure reduction. An expansion-type CNC is typically a batch instrument. Expansion-type CNCs have been used in atmospheric aerosol research for many years. A continuous CNC (for example, as disclosed by Agarwal and Sem, 1980) is widely used. In the continuous CNC, supersaturation is formed by cooling a laminar aerosol flow that is saturated with working fluid vapor. In the conductive cooling type continuous CNC the aerosol-containing sample is typically saturated with working fluid and then cooled whereby the working fluid condenses on the aerosol particles. However, one disadvantage of the conductive cooling CNC is its sensitivity to moisture in the sample gas. Moisture in the sample gas may also condense on the aerosol particles or otherwise interfere with the condensation of the working fluid and affect the measured particle count. Attempts to remove moisture from the sample gas typically can also remove aerosol particles, which hampers the accuracy of the particle measurements. Aspects of the present invention overcome this deficiency of the prior art.\nThe third type of CNC, known as a turbulent mixing CNC (or TMCNC), is based on turbulent mixing of a gas flow with particles with working fluid vapor. The TMCNC instrument has not yet been commercialized. One prior art TMCNC is described by Kogan and Burnashova (1960) and was further developed in different versions by Okuyama, et al. (1984); Ankilov, et al. (1991); Kousaka (1993); and Mavliev and Wang (1999). The major advantage of the TMCNC is the flexibility of generating supersaturation by simply mixing the aerosol flow and a separate gas flow saturated with the working fluid vapor.\nThe CNC is one of the most sensitive devices for detecting nanometer-size particles, for example, some CNC studies have reached a detection limit of 2-3 nanometers (nm) (Stolzenburg and McMurry, 1991; McDermont et al., 1991; Okuyama et al., 1984; Mavliev and Wang, 1999). In one commercially available CNC system that can detect 3 nm particles, the aerosol stream is directed through a capillary tube in order to align the particle stream within a small central zone of uniform saturation conditions. However, the capillary in this system is prone to problematic clogging when used to directly measure ambient aerosols and cleaning the capillary tube can be difficult.\nThe minimum detection efficiency of CNCs is very sensitive to the size of particles being detected (Makela, et al., 1996). The detection efficiency may also be sensitive to the composition of the particle (Mavliev et al., 2001) and to the nature of the working fluid (Lee et al., 2003; Mavliev et al., 2004).\nAlthough the CNC is primarily devoted to measuring the number concentration of particles, some prior art CNCs have been shown to be able to measure size distribution of particles, for example, of nanometer-sized particles. In one prior art system, the size distribution of nuclei can be measured by means of changing the CNC's sensitivity (McDermont et al. 1991) and by means of measuring the size of grown particles (Ahn and Liu, 1990; Rebours, et al., 1996; Saros, et al., 1996). One prior art method is based on the fact that the growth of smaller particles is delayed because of the Kelvin effect that results in the final size of particles being dependent on initial nucleus size. However, these prior art devices provide unsatisfactory size resolution and their operational range is limited to a range of 3-10 nm. In addition, in some of these prior art systems, the growth time for particles of the same size depends strongly on spatial uniformity of supersaturation. For most continuous flow CNCs, the spatial distribution of supersaturation is not uniform because of the use of diffusive cooling in the laminar flow (Stolzenburg, 1988).\nMoreover, some commercially available prior art CNCs were developed as laboratory research tools or to monitor clean rooms, and, as far as the present inventors know, no effort has yet been made to optimize their performance for ambient aerosol monitoring. For example, typically, the aerosol-laden air streams of such research CNCs are not conditioned in any way, nor is there any control or limitation of the introduction of large particles to these CNCs. In prior art ultrafine CNCs, large particles can clog the device, for example, the capillaries in such devices. In addition, in prior art systems there is typically no control of temperature and humidity in the inlet stream and both temperature and humidity can produce some variability in the instrumental response. These deficiencies in the prior art devices cause problems when using a CNC for ambient particle monitoring.\nIn addition, prior art CNCs typically use n-butanol as the working fluid. However, butanol is toxic, flammable, has a noxious odor, and is thus undesirable in a commercially-available device that may readably experience human contact, operation, and servicing.\nThus, in view of the deficiencies and disadvantages of prior art CNCs, there is a need for an improved stand-alone TMCNC aerosol particle counter that will provide improved sensitivity, for example, for particle sizes down to the order of about 2 nm, well-defined inlet characteristics, robust design, consistent and reproducible performance, and employ a working fluid that is less offensive and less dangerous to humans. Aspects of the present invention address these and other deficiencies of the prior art while providing improved detection efficiency, greater ease of use, and adaptability for use in ambient gas particulate measurements."}
{"text": "Conventional optical data storage devices are configured to read digital data from and write digital data to a removable optical disc. Currently, writable compact discs (CD-R) and re-writable compact discs (CD-RW) are popular formats for personal computers and other like devices. Re-writable digital versatile discs (DVDs), known as DVD-RAMs (random access memory), DVD-R, DVD-R/W, etc., have also become popular as the price of the applicable DVD drive devices has become more affordable.\nThe process of writing data to an optical disc is often referred to as “burning” the disc, since a beam from a write laser is used to selectively raise the temperature of certain materials within the optical disc such that the materials are altered in some manner. Consequently, features are formed on the disc. These features represent binary data values, i.e., 1's and 0's, which can subsequently be detected (read) using a read laser. After burning in a CD-R, for example, the user may decide to label the CD-R disc to reflect the data files that have been stored on the disc.\nThe LightScribe™ printing system developed by the Hewlett-Packard Company of Palo Alto, Calif. is a direct disc labeling technology that produces optically visible text and graphics on the surface of a CD or DVD disc media using the disc drive's laser and a specially coated LightScribe disc. A LightScribe disc drive writes data (on the data side of the disc) and creates optically visible labels (on the label side of the disc) with the same drive, and may use the same laser for both operations by flipping the disc over after writing the digital data in order to create the label. Using laser technology, LightScribe's specially coated discs chemically change when a laser strikes its label side with energy of a particular wavelength, producing optically visible marks that can be arranged to form titles and graphics using LightScribe software. One embodiment of a suitable coating is described in U.S. Patent Application Publication No. 2003/0108708 A1.\nIt would be useful to avoid having to flip a disc over in order to label it. It would also be useful to reduce the cost or complexity of a printhead for a disc drive."}
{"text": "1. Field of the Invention\nThe present invention relates to ammunition, especially for utilization with cannon which are mounted on board aircraft, including a subcaliber projectile core and a propulsion mechanism.\n2. Discussion of the Prior Art\nAmmunition of that type is described in the specification of German Patent No. 31 31 540. A propulsion mechanism which is constructed in the form of a sabot is interconnected with the subcaliber-sized projectile core through the utilization of a locking or spring ring which is provided to act as a securing device under centrifugal force. In response to the spin acting on the projectile, the spring ring will expand so as to facilitate the projectile core exiting from the sabot. The sabot possesses a high flow resistance so as to fall back during flight relative to the projectile core. Consequently, this will produce the danger of the sabot, as well as the plastic components of the projectile hood, striking against and damaging the aircraft which is firing the ammunition.\nIn the disclosure of German Laid-Open Patent Appln. No. 35 08 053 there is described a subcaliber-sized projectile possessing a sabot which is constituted from a plurality of segments. The sabot is dispersed or split up into its segments as a consequence of its high flow resistance, so as to cause the projectile core to be freed. When a projectile of this kind is fired from a flying aircraft, there is then encountered the danger that the aircraft will fly into the segments from the sabot."}
{"text": "Software-defined networking (SDN) is an approach to networking in which control is decoupled from hardware and given to a software application, referred to herein as a SDN controller. Conventionally, when a packet arrives at a network device in a typical network, the control plane, which is based on the switch's proprietary firmware, tells the device where to forward the packet. In some network devices, the control plane may be configured to recognize different types of packets and treat them differently, such as forwarding the packets to different network devices or dropping the packets if certain criteria are met.\nIn a SDN enabled network, a network administrator can change how the network devices route data packets using the software application without having to, for example, send a technician to the individual network devices. The administrator can remotely reconfigure the rules in the control plane of the network devices—e.g., prioritizing, de-prioritizing or even blocking specific types of packets—using the SDN controller. A SDN enabled network may allow the administrator to manage traffic loads in a flexible and more efficient manner relative to conventional network devices that use control planes based on firmware.\nMoreover, SDN allows a network administrator to configure a switching fabric across multi-vendor hardware and different proprietary firmware. One standard that permits a network administrator to remotely configure and control the control plane of network devices is the OpenFlow® standard (OpenFlow is a trademark of the Open Networking Foundation). Generally, OpenFlow lets network administrators control routing tables for a plurality of network devices using a single SDN controller.\nTo facilitate understanding, identical reference numerals have been used, where possible, to designate identical elements that are common to the figures. It is contemplated that elements disclosed in one embodiment may be beneficially utilized on other embodiments without specific recitation.\nDescription\nOverview\nEmbodiments described herein include a method, a system and a computer program product for selecting a tag for use in a SDN environment. The method, system, and computer program product receive an attribute from a user application executing on a computing device, the attribute describing a network flow associated with the user application. The method, system, and computer program product comparing the attribute to one or more rules to select at least one tag for the network flow, where the selected tag is associated with a corresponding action to be performed in the SDN environment. The method, system, and computer program product insert the received tag into at least one packet in the network flow. After inserting the selected tag, the method, system, and computer program product transmit the at least one packet into the SDN environment, where at least one network device in the SDN environment is configured to recognize the tag in the at least one packet and perform the corresponding action.\nDescription of Example Embodiments\nSoftware-defined networking (SDN), in general, provides additional flexibility and solidarity relative to conventional networks. Using an SDN controller, which may be either centrally located or located on the respective network devices, a network administrator can configure the control plane and dictate how the network device routes data. For example, the network administrator may assign criteria or rules that, when satisfied, instruct the network device to perform a specific action on the received packet—e.g., drop the packet, forward the packet to a particular network device, evaluate the packet using an application on the network device, and the like. In one embodiment, the SDN controller configures the routing table or forwarding table (i.e., forwarding information base) in a network device based on the network administrator's preferences.\nAfter the packet has been generated by the user application, a network device in the SDN enabled network may evaluate the packet to determine what action to take. The network devices may rely on the different headers of the packet—e.g., source and destination addresses, security group tag, MAC-layer priority value, VLAN ID, etc.—to determine the action. Alternatively or additionally, the network device may perform a deep packet inspection to determine what action is appropriate. In one embodiment, the network administrator may prefer that the network devices only have to evaluate the packet headers or other superficial information to identify the action to perform which may require less time than a deep packet inspection. However, the headers may not be able to convey information to the network device that is useful when determining the action to perform such as the duration of the flow, the amount of data that will be transferred in the flow, certain contractual requirements related to the flow, and the like. Moreover, in some instances, even performing a deep packet inspection may not reveal the desired action. For example, if the network administrator guaranteed that a certain network flow generated by a user application will be routed on the least-latency path through the network, evaluating the headers or performing a deep packet inspection on the packet may not reveal the source application, and thus, the network device may not be able to identify the packets that should be routed along the least-latency path.\nInstead, the user application may provide a hint which may be used by a tag controller to identify a tag that is inserted into the packets generated by the user application. In one embodiment, when the user application is executed on a computer device, the application transmits information (i.e., a hint) about its associated network flow via a stack associated with the computing device. One of the layers of the stack—e.g., the operating system or a virtual machine in the computing device—may be communicatively coupled to the tag controller. This layer may receive from the tag controller a hint database that includes information for selecting a tag based on a hint received from the user application. For example, the hint may identify the user application as a high-speed trading application that requires the path through the network with the least latency. Using the hint database, the layer may tag each packet generated by the trading application to a particular VLAN. Using an SDN controller, the network administrator may configure each network device such that the packets with the assigned VLAN tag traverse the network using the least latency path. In this manner, the network administrator uses the hint provided by the user application to select a tag which may be inserted into the packet which the network devices can easily evaluate to identify the appropriate action.\nBecause the tag controller provides the hint database to one of the layers in the stack, the network administrator may not need to directly monitor the user application. Instead, the layers in the stack may select the appropriate tag using the hint database and insert the tag into the network flow. Similarly, the programmer who developed the user application may not need to have knowledge about the underlying layout of the SDN enabled network. Instead, the programmer of the user application may provide a description of the hints to the network administrator which may then load the hint database onto the tag controller. The controller may then provide the database to one of the layers in the stack. Thus, the hint and the hint database are an indirect communication method by which the stack can select the appropriate tag for the network traffic that suits the configuration of the particular network.\nFIG. 1 illustrates a system 100 for using information provided by a network traffic application to generate a tag for a SDN enabled network. System 100 includes a tag controller 105, a computing device 110 that hosts a network traffic application 120 (e.g., a user application), a SDN controller 125, an edge device 130, and a SDN enabled network 135. The tag controller 105 may be an application executing on a computing device separate from the computing device 110 (e.g., a server connected to network 135). In one embodiment, the tag controller 105 may be accessible to a network administrator of network 135 while the network traffic application 120 is not. Accordingly, to aid the administrator and tag controller 105 when selecting an appropriate tag for the network flow (or flows) generated by the network traffic application 120, the tag controller 105 may transmit a hint database (not shown) to computing device 110 via the communication path 140. The user application 120 may provide information about the network flow—i.e., a hint—which the computing compares to the hint database to select the appropriate tag. In another embodiment, the computing device 110 may use communication path 140 to transmit the hint to tag controller 105. In this scenario, tag controller 105 may select the appropriate tag and transmit the tag back to computing device 110. The hint may define the duration of the network flow, the number of bytes likely to be transferred by the network flow, a description of the end points of the flow, an identity of the network traffic application 120, contractual obligations associated with the flow, and the like. In one embodiment, the network traffic application 120 may provide the hint before application 120 begins to generate network traffic—e.g., application 120 provides the hint when it is first executed in the computing device 110. Network traffic application 120 may provide a hint based on other predefined events, such as opening a remote connection, estimating a duration of a network flow or the amount of data a network flow will carry, receiving a query or request for data, and the like.\nCommunication path 140 may be any communication channel used to couple the computing device to the tag controller 105. For example, communication path 140 may include edge device 130 and network 135 which route the hint from the computing device 110 to the tag controller 105. Alternatively, computing device 110 may have a separate communication link that permits device 110 to directly receive the hint database from the tag controller 105.\nComputing device 110 may be any computing structure such as a server, personal computer, database, and the like that uses on a processor 115 to execute one or more applications. Processor 115 represents any processor capable of performing the task discussed herein and may include multiple processors or processors with multiple cores. When executed by processor 115, the network traffic application 120 generates one or more network flows that traverse edge device 130 and the SDN enabled network 135 to one or more end points (not shown). That is, the network traffic application may be loaded into a memory associated with the computing device 110. Network traffic application 120 may be any type of user application, such as a web server application, database application, remote-storage application (e.g., cloud storage application), and the like.\nAs shown, computing device 110 is coupled to edge device 130 (e.g., a switch, router, bridge, or other network device) that provides access to network 135. The SDN enabled network 135 may include a plurality of network devices that route the network flow to one or more end points. In one embodiment, the SDN enabled network 135 includes network devices whose control planes are configurable by the SDN controller 125. Specifically, the network administrator may reconfigure the rules of the control planes using the SDN controller 125 in order to prioritize, de-prioritize or block specific types of packets.\nIn one embodiment, the SDN enabled network 135 may represent only the network devices whose control planes are configurable by the SDN controller 125. That is, network 135 may be connected to other network devices that are controlled by a different SDN application or network devices that are not SDN enabled. For example network 135 may be only a portion of a larger network, e.g., the Internet. In these scenarios, the network administrator may have permission to access and configure only a subset of the network devices used to transmit the network flow.\nAs shown, tag controller 105 and SDN controller 125 may be communicatively linked via path 145. Like path 140, communication path 145 may use the SDN enabled network 135 to transmit data between the controllers 105 and 125. Alternatively, controllers 105 and 125 may be applications executing on the same computing device. Generally, communication path 145 permits the controllers 105 and 125 to share information about network 135. For example, SDN controller 125 may provide information to the tag controller 105 about the particular configuration of the network 135. Based on this information, the tag controller 105 may update the hint database such that the hints provided by application 120 are matched to tags customized to the configuration of the SDN enabled network 135. In one embodiment, the SDN controller 125 and the SDN enabled network 135 may follow the OpenFlow® standard for configuring the rules used to route data in network 135.\nFIG. 2 illustrates a method 200 for generating a tag for a SDN enabled network using information about a network traffic application. At block 205, the network traffic application generates information about a network flow (i.e., a hint). The hint may be either a standalone data chunk or be embedded within another data message or communication. In one embodiment, the hint may be separate from the network flow. For example, the hint may be transmitted from the network traffic application to an operating system executing on the host computing device as part of loading the network traffic application into memory.\nAt block 210, the computing device selects a tag based on the hint issued by the network traffic application. For example, the operating system on the computing device may have a hint database that includes a plurality of rules that use the hint as an input to select the appropriate tag. Moreover, in one embodiment the computing device may also consider information in the packets when selecting the tag. An example rule could be if TCP flow is to port 80 (which is gathered by evaluating packets of the network flow) and transfer size is less than five megabytes (which may be information gathered from the hint) then the appropriate tag is a MPLS label value of 2459. The different manner by which the hint may be transmitted from the network traffic application to the tag controller and how the hint database is transmitted from the tag controller to the computer device will be discussed in further detail below. Once the hint is received, the computing device evaluates the hint using the hint database to determine the appropriate tag for the network flow or flows associated with the hint. As mentioned previously, the hint may define the duration of the network flow, the number of bytes likely to be transferred by the network flow, a description of the end points of the flow, an identity of the network traffic application, contractual obligations associated with the flow, and the like. In one embodiment, the network traffic application may transmit the hint before generating any network traffic on the network flow. Alternatively, the network traffic application may provide the hint after starting to transmit packets in a network flow on the SDN enabled network.\nIn one embodiment, the computing device may transmit the hint to the tag controller. Using the tag controller, the network administrator may select the appropriate tag for a network flow. For example, the hint database may not include an appropriate tag for a received hint. Instead of not selecting a tag, the computing device may transmit the hint to the tag controller. The administrator may use a display screen communicatively linked to the tag controller to view the hint received from the network traffic application 120. For example, the hint may inform the network administrator that the network flow transmits high-priority information that needs to traverse the network quickly. In response, the network administrator may select a tag that assigns the packets in the network flow a high-priority. In this manner, the tag controller may select a tag using input from the network administrator. Moreover, the network administrator may then use the SDN controller to configure the SDN enabled network to route packets with the selected tag through the least-latency path.\nIn one embodiment, the network traffic application may generate a plurality of different network flows. Each of these network flows may be associated with a different hint and/or tag. In one embodiment, the network traffic application may provide a hint associated with each of the flows, and at block 210, the computing device may evaluate each of the hints independently and select a tag for each of the flows. The selected tags may be the same or different. For example, the network traffic application may generate two different network flows: one flow is dedicated to transmitting confidential data while the other flow is associated with data that is time sensitive. The computing device may assign two separate tags to the flows such that first flow uses a firewall when traversing the network while the second flow uses the least-latency path through the network. In another embodiment, the network traffic application may provide only one hint even if the application generates a plurality of different flows. Nonetheless, the computing device may use the single hint to issue multiple tags. For example, if the hint identifies the network traffic application as a web server, the hint database in the computing device may be preconfigured with multiple rules that consider attributes of the network flows (e.g., which ports the flows use). Using the rules, the computing device may select a tag for each of the network flows from the web server.\nAt block 215, the computing device or the network traffic application may insert the selected tag into the network flow. For example, the computing device or network traffic application may selectively insert the tag into only certain packets in the network flow or insert the tag into all the packets in the flow. Generally, the tag may be any value that, when inserted into a packet, is easily identifiable by a network device—i.e., a value that is found without performing a deep packet inspection. In one embodiment, the tag is inserted into a packet header. In this scenario, the tag may be a predefined value that allows a network device to distinguish the packet, and thus the network flow, from other packets and network flows. For example, the tag may be a particular value for a packet's multiprotocol label switching (MPLS) header, a security group tag, a service insertion tag, a differentiated services code point (DSCP) value, VLAN ID, and the like. As will be discussed below, the tag may be inserted into the packets of the network flow by different elements in the computing device—e.g., the tags may be inserted by the network traffic application, the operating system, or a virtual machine executing on the computing device.\nAt block 220, the network administrator may configure the SDN enabled network to perform an action based on identifying the tag in a packet received at a network device. Stated differently, the tag may cause a network device to perform a certain action. This action may include routing the packet to a particular network device, processing the packet using an application executing on the network device, manipulating the packet, discarding the packet, and the like. In one embodiment, the network device may include a table (e.g., a flow table) that matches the tag with the action the network device should perform. In one embodiment, the network device may remove the tag from the packet before forwarding the packet to a different network device.\nIn one embodiment, the network administrator may configure the network in response to the tag assigned to the network flow. For example, if the hint indicates that the network flow should pass through a firewall when traversing the SDN enabled network, the computing device may assign a tag that routes the packets of the network flow through a network device hosting the firewall. In turn, the network administrator may use a SDN controller to access the edge device connected to the computer device hosting the network traffic application. Specifically, the network administrator may change or update the routing rules (or control plane) of the edge device such that packets containing the tag are routed to the network device hosting the firewall. As the network traffic application begins to generate packets, the tags are inserted into the packets which are detected by the edge device. This device then routes the packets through the firewall. Thus, in this example, the network administrator may configure the SDN enabled network after receiving the hint from the user application—e.g., the network traffic application.\nIn another embodiment, the tag controller or the network administrator may generate the hint database based on the current configuration of the SDN enabled network. For example, the network may have been previously configured (i.e., before the hint was transmitted at block 210) such that all data packets assigned to a particular VLAN will use the least-latency data path to traverse the network. Based on this knowledge, the tag controller may put a rule in the hint database indicated that the network flow should be assigned the particular VLAN ID associated with the least-latency data path. Accordingly, in this example, the configuration of the SDN enabled network is considered before assigning the tag to the network flow or receiving the hint from the network traffic application. Furthermore, the network administrator may not reconfigure the SDN enabled network based on the received hint, but instead uses the current configuration of the network to select the appropriate tag for the network flow.\nRegardless of whether the network administrator changes the network configuration after the application begins to transmit hints or not, the configuration of the network may nonetheless be used to suit the needs of the particular user application without the application's developer having any knowledge of the configuration of the network. Accordingly, the network traffic application does not need to be customized for a particular network configuration. Instead, by transmitting the hint, any SDN enabled network may provide a customized tag that causes the underlying network devices to perform an action corresponding to the particular configuration of the network.\nAlthough method 200 discusses selecting a tag based on a hint associated with a particular network flow (e.g., a plurality of related data packets), the present disclosure is not limited to such. For example, in other embodiments, the computer device may select a tag on a per-packet basis. That is, the network traffic application may transmit a hint which the computer device uses to identify a tag for each packet individually.\nFIG. 3 illustrates a system 300 for using information about a network application to generate a SDN tag using a hierarchical stack. Specifically, system 300 includes one or more computing devices 110 that may host one or more virtual machines 305. In one embodiment, virtual machine 305 emulates the computer architecture and functions of a physical computing device. That is, virtual machine 305 may be a software implementation of a computing device that executes programs and applications like a physical machine. In this manner, the hardware resources of the computing device 110 may be divided into a plurality of virtual machines 305. Each virtual machine 305 may include one or more instantiations of an operating system 310 which may be used to execute the network traffic application 120. Operating system 310 may be any operating system capable of performing the functions described herein.\nNetwork traffic application 120 includes a network library 320 which may contain the library routines for the application 120. In one embodiment, the network library 320 may be a sockets library. When the network traffic application 120 first is loaded into memory, opens a disk file, opens a remote connection, or registers to accept incoming connections, the application 120 may supply a chunk of data (i.e., the hint) to another element in the computing device 110. As mentioned previously, the hint may be a connection type, an estimated length of connection in both total data transfer and time, a characterization of the bandwidth required, and an indication of the reliability of these estimates. For example, when opening a TCP connection, the hint may be transmitted to the network library 320 which then passes the hint to operating system 310, but when opening a disk file, the hint may provided directly to operating system 310, thereby skipping network library 320.\nThe other elements shown in FIG. 3—e.g., controllers 105 and 125, network 135, and edge device 130—may perform a similar function as that described in FIG. 1. However, in one embodiment, edge device 130 may be virtualized and hosted within computing device 110.\nFIG. 4 illustrates a system 400 that uses information about the network application to generate a SDN tag. Specifically, system 400 includes a stack 405 through which the network flow and the hint generated by the network traffic application 120 may traverse. Specifically, arrow 410 illustrates the path the hint may travel in stack 405 while arrow 415 illustrates the flow of network traffic from the application 120 to the SDN enabled network 135.\nDotted lines 420A-C illustrate possible communication links between the different layers in stack 405 and the tag controller 105. In one embodiment, at least one of the layers of stack 405 is communicatively coupled to the tag controller 105. For example, the operating system 310 may be configured to communicate with the tag controller 105 (i.e., path 420B is used) while the network library 320 and the virtual machine 305 are not. However, in other embodiments multiple layers of stack 405 may be communicatively coupled to tag controller 105 via lines 420A-C. Although lines 420A-C are illustrated as being separate from network 135, in one embodiment, the network library 320, operating system 310, or virtual machine 305 may communicate with tag controller 105 via the network 135. That is, the layer of stack 405 that is configured to communicate with the tag controller 105 may receive the hint database from the controller 105 through the SDN enabled network 135 and edge device 130. Alternatively, the computing device may provide a separate communication link for transmitting the hint database from tag controller 105 one of the layers in stack 405.\nIn one embodiment, the hint may be a data chunk ranging from, for example, one or more bits to twenty bytes of information. In one embodiment, the hint may be transmitted using I/O control in the virtual machine 305. The hint may be embedded into other data communications transferred between the layers in stack 405. For example, when detecting a remote connection, the application 120 may embed the hint in the I/O control data that flows between the different layers of stack 405. Even if one of the layers is not configured to process the hint (i.e., the layer does not have a hint database), the layer may pass the embedded hint to a lower layer which does have the hint database. For example, assuming that the operating system 310 is not configured to communicate with the tag controller 105 to receive the hint database but the virtual machine 305 is, once the operating system 310 receives the hint, the operating system 310 may simply pass the hint to the lower layer—i.e., the virtual machine 305—which recognizes the hint and selects the appropriate tag using the database. In this manner, the different layers of stack 405 do not need to be configured to recognize or process the hint.\nAlthough path 410 illustrates that the hint traverses through the network library 320, operating system 310, and virtual machine 305, this is not a requirement—i.e., the hint may pass through only one or two of the layers. For example, if the network library 320 is configured to process the hint, then the hint may not need to be passed on to the lower levels of stack 405. If the hint is embedded in other data communication (e.g., I/O control), then the hint may “piggy-back” on that data until the hint reaches the layer with a hint database. That layer may, if desired, remove the hint from the carrier data communication before passing the data to a lower level. Thus, in one embodiment, the hint may traverse stack 405 until the hint reaches the layer that selects the tag.\nMoreover, the hint may skip over one or more of the layers 405 when traversing the stack 405. For example, in some instances (e.g., when opening a disk file), the application 120 may provide the hint directly to operating system 310 and skip network library 320.\nSystem 400 also illustrates that a network administrator 425 may influence which tag is selected as well as configure the network 135 using SDN controller 125. As mentioned previously, the network administrator 425 may use tag controller 105 to select a particular tag or re-configure the control planes of the network devices in network 135 to perform a certain action when detecting the tag in a packet. However, in other embodiments, the techniques discussed herein may be performed without any input from the network administrator 425.\nFIG. 5 illustrates a method 500 for using information about a network application to generate a SDN tag using a hierarchical stack. At block 505, the tag controller may receive a description of the hints used by a network traffic application. In one embodiment, the programmer who develops the network traffic application may publish a description of the various hints the application may provide. A network administrator may access the list and upload it to the tag controller. Alternatively, the tag controller may periodically access, e.g., a central data store to retrieve descriptions of hints user applications may provide. These descriptions may be used to decode a hint transmitted by the network traffic controller to the tag controller. More specifically, the network administrator or tag controller may use the hint descriptions to generate the rules for the hint database. For example, the network traffic application may use a predefined bit value as a hint to indicate that the network flow is expected to be less than a gigabyte (e.g., “0”), while the bit value “1” indicates that the data flow is expected to be greater than a gigabyte. The tag controller may use the description to decode and interpret the hint received from the application to generate one or more rules for assigning a tag based on this hint. In this manner, a programmer may use any code when transmitting a hint and then provide the description of the code to the tag controller using a different communication method—e.g., publishing the description of the hint on a data repository that is accessible to the tag controller or the network administrator. Doing so may reduce the size of the data hint—e.g., the hint may range from one bit to twenty bytes. Once the hint descriptions are generated into rules and placed in the hint database, the tag controller may transmit the hint database to one of the layers in the stack.\nIn one embodiment, the tag controller may receive the description of the hint before receiving the hint from the network traffic application. For example, the tag controller may periodically download the description of the hints from a central repository or the network administrator may selectively load descriptions of hints into the tag controller based on the type of user applications the network administrator knows will use the SDN enabled network. In one embodiment, the tag controller of network administrator may retrieve the description of the hints after receiving the hint from the network traffic application. The tag controller may receive a coded hint, and if the controller cannot match the hint to a rule in the hint database, the tag controller may attempt to download the description from, for example, the central repository and transmit an updated hint database to one of the layers in the stack which can decode the hint and select the appropriate tag.\nIn other embodiments, however, the hint may not need a description. That is, the hint provided by the network traffic application may include a description that explains what the hint is. For example, the hint may be the identity of a user that opened a remote connection or the estimated duration of a network flow. In these scenarios, the hint may be larger than the coded hints discussed above but this may enable method 500 to skip block 505. In other words, the hint may be designed such that the computing device can use the hint to select a tag using, for example, a generic version of the hint database without receiving any supplemental description of the hint.\nAt block 510, the network traffic application may issue a hint associated with one or more network flows generated by the application. As shown in FIG. 4, the network traffic application may issue the hint along a similar data path (e.g., stack 405) as the packets associated with the network flow. In one embodiment, the network traffic application may determine the hint to issue based on an event occurring at the network traffic application. That is, the network traffic application may issue a plurality of different hints that each correspond to a different event. Non-limiting examples of these events include the application first being executed, opening a remote connection, estimating a duration of a network flow or the amount of data a network flow will carry, receiving a query or request for data, and the like. In response to the event, the network traffic application may select a particular hint (or hints). For example, if network traffic application receives a request to transmit a particular table in a database to a user, the application may issue one hint that identifies the user making the request and a second hint providing an estimated duration of the network flow in order to satisfy the query. The network traffic application may issue these hints simultaneously or sequentially.\nAt block 515, the hint issued by the network traffic application passes through at least one layer in the stack 405 shown in FIG. 4. In one embodiment, the hint or hints may pass through the stack sequentially—e.g., through the network library (e.g., a socket library), operating system, and virtual environment. Alternatively, the hint may pass through only a subset of these layers or may skip layers as the hint traverses the stack. Further still, the hint may pass through the stack only until the hint reaches the layer that is communicatively linked to the tag controller and contains the hint database. If, for example, the hint is embedded in a different data communication, the layer that selects the tag may remove the hint before passing the data communication to a lower level in the stack. In one embodiment, even if a layer of the stack is not tasked with transmitting the hint to the tag controller, the layer may change or amend the hint before passing the hint to a lower level.\nIn one embodiment, one of the layers in the stack (e.g., a virtual switch or the edge device) may have a flow table (e.g., a superset of the OpenFlow flow table) which may include mask-and-match rules for the SDN enabled network. Specifically, the flow tables may identify a tag or marker that indicates to a SDN enabled network device to take a corresponding action. Actions may include forwarding packets with the tag or marker out a certain port, adding or removing an encapsulation, forwarding through a tunnel to another part of the network, and the like. The SDN controller may publish these rules on the various network devices in the SDN enabled network or even on a virtual switch located in computing device 110 shown in FIG. 1.\nAt block 520, the layer of the stack with the hint database selects the appropriate tag for the network flow. For example, the operating system on the computing device may identify the hint and select the tag. As mentioned previously, the hint may define the duration of the network flow, the number of bytes likely to be transferred by the network flow, a description of the end points of the flow, an identity of the network traffic application, contractual obligations associated with the flow, and the like. In one embodiment, the network traffic application transmits the hint before generating any network traffic. Alternatively, the network traffic application may provide the hint after starting to transmit a network flow on the SDN enabled network.\nIn one embodiment, the layer with the hint database may transmit the hint to the tag controller. This may occur, for example, if the hint database does not have a rule matching the hint. Alternatively, instead of storing a local copy of the hint database in the layers of the computing device, the layer may transmit each hint to the tag controller which selects the tag using the database. The tag controller may then transmit the tag back to the layer which can insert the tag into the network flow.\nIn one embodiment, the network traffic application may generate a plurality of different network flows. Each of these network flows may be associated with a different hint and/or tag. Or the network traffic application may generate a plurality of hints for a single network flow. At block 520, the computing device may evaluate each of the hints independently and select a tag for each of the flows. These tags may be the same or different. In another embodiment, the network traffic application may provide only one hint even if the application generates a plurality of different flows. Nonetheless, the computing device may use the single hint to select multiple tags or select multiple tags for a network flow. Furthermore, if the network traffic application issues multiples hints for a single network flow, the computing device may select a tag for each of the hints.\nAt block 525, the layer that selects the tag may insert the tag into the network flow. For example, the layer may selectively insert the tag (or tags) into a subset of the packets in the network flow or insert the tag into all the packets in the associated network flow. As mentioned previously, the tag may be any value that, when inserted into a packet, is easily identifiable by a network device—i.e., a value that network device can identify without performing a deep packet inspection.\nAt block 530, the network administrator may configure the SDN enabled network to perform an action based on identifying the tag in a packet received at a network device. Stated differently, the tag may cause a network device to perform a certain action. This action may include routing the packet to a particular network device, processing the packet using an application executing on the network device, manipulating the packet, discarding the packet, and the like. In one embodiment, the network device may include a table that matches the tag to the corresponding action the network device should perform. In other embodiments, the network device may remove the tag from the packet before forwarding the packet to a different network device.\nIn one embodiment, the network administrator may configure the network in response to the tag assigned to the network flow using, for example, the SDN controller whose function was discussed previously. Specifically, the network administrator may change or update the routing rules (or control plane) of the edge device such that packets which include the tag are routed to a network device hosting the firewall. As the network traffic application begins to generate packets, the tags are inserted into the packets and the edge device detects the tags and routes the packets accordingly. Thus, in this example, the network administrator may configure the SDN enabled network after receiving the hint from the user application—e.g., the network traffic application.\nIn another embodiment, the tag controller or the network administrator may base the tag selection on the current configuration of the SDN enabled network. That is, the network administrator may not reconfigure the SDN enabled network based on the received hint, but instead uses the current configuration of the network to select and appropriate tag and action for the network flow.\nRegardless of whether the network administrator changes the network configuration after the hint is received or not, the configuration of the network may nonetheless be used to suit the needs of the particular user application without the application's developer having any knowledge of the configuration of network. Accordingly, the network traffic application does not need to be customized for a particular network configuration. Instead, by transmitting the hint, any SDN enabled network may provide a customized tag that causes the underlying network devices to perform an action corresponding to the particular configuration of the network.\nAlthough method 500 discusses selecting a tag based on a hint for a particular network flow (e.g., a plurality of related data packets), the present disclosure is not limited to such. For example, in other embodiments, the computing device may select a tag on a per-packet basis. That is, the network traffic application may transmit a hint which the computing device uses to identify a tag for each packet individually.\nFIGS. 6A-6C illustrate initiating a stack for selecting a tag based on received hints, according to embodiments described herein. That is, FIGS. 6A-6C may occur before a hint is issued by a user application. Specifically, FIG. 6A illustrates a system before a virtual machine is initialized in the computing device 110. The computing device 110 includes a hypervisor 605 which may be a portion of software or firmware that creates and manages virtual machines. Here, hypervisor 605 includes a virtual switch 610, but in other embodiments, the virtual switch 610 may not used. The hypervisor 605 may establish communication links 615 and 620 with tag controller 105 and SDN controller 125, respectively. As will be discussed below, communication link 615 may enable the hypervisor 605 to establish a communication link between the tag control 105 and other elements in the computing device 110.\nThe SDN controller 125 may also establish a communication link 625 to edge device 130 as well as the communication link 620 to the virtual switch 610. Through these links, the SDN controller may establish SDN rules for interpreting tags inserted in received packets. Although two links are shown, in one embodiment, the SDN controller 125 may be communicatively coupled to only one of the links.\nIn FIG. 6B, hypervisor 605 loads virtual machine 305 into the memory of computing device 110 and begins to execute the operating system 310. Using communication link 615, hypervisor 605 may inform tag controller 105 about virtual machine 305 (or vice versa). Accordingly, tag controller 105 or virtual machine 305 may establish communication link 630 which enables tag controller 105 to pass data (e.g., a hint database) to the virtual machine 305 and/or the operating system 310.\nIn one embodiment, tag controller 105 may use communication link 615 to load a hint database 635A into the hypervisor 605. As discussed previously, the hint database 635A may be located in multiple layers in the stack or at only select layers and may be used by the layer to select a tag to insert into one or more packets of a network flow based on a hint provided by a user application (e.g., the network traffic application).\nIn addition, the SDN controller 125 may use communication links 620 and 625 to load flow tables into 640A-B. In one embodiment, SDN controller 125 may load the flow table 640 into one of the networking elements instead of both as shown here. The flow tables 640A-B may include mask-and-match rules (e.g., SDN rules) for the SDN enabled network. Specifically, the flow tables may identify a tag or marker that indicates to a SDN enabled network device to take a corresponding action. Thus, if the layers above the virtual switch 610 and edge device 130 insert a tag into a packet, the flow tables 640A-B enable the network elements to identify the tags and determine the corresponding action to perform.\nIn FIG. 6C, operating system 310 may load the network traffic application 120 into memory. Here, the operating system 310 may provide application information to tag controller 105. In one embodiment, this application information may be the hint discussed previously. To interpret the hint, tag controller 105 may provide the operating system 310 or virtual machine 305 with a hint database 635B. Moreover, in one embodiment, the tag controller 105 may update the hint databases 635A-B if the tag controller 105 learns more about the hints issued by the network traffic application 120.\nIn one embodiment, network traffic application 120 may continue to transmit hints to the elements in computing device 110. That is, instead of receiving only one hint from the network traffic application 120 once the application is first loaded into memory, application 120 may issue new hints based on other events such as establishing a new network connection, receiving a user request, or estimating a duration or total data transfer amount of the network flow.\nTag controller 105 may also establish a communication link 632 with the network traffic application 120. Although not shown, tag controller 105 may use this link 632 to load a local copy of the hint database into the network library 320 as discussed above. In this manner, any of the elements (or layers) in computing device 110 (e.g., virtual machine 305, operating system 310, network library 320, or hypervisor 605) may be configured to be in communication with tag controller 105. Using these links, the tag controller 105 may provide updates the hint databases 635.\nAs will be appreciated by one skilled in the art, embodiments may be embodied as a system, method or computer program product. Accordingly, aspects may take the form of an entirely hardware embodiment, an entirely software embodiment (including firmware, resident software, micro-code, etc.) or an embodiment combining software and hardware aspects that may all generally be referred to herein as a “circuit,” “module” or “system.” Furthermore, aspects may take the form of a computer program product embodied in one or more computer readable medium(s) having computer readable program code embodied thereon.\nAny combination of one or more computer readable medium(s) may be utilized. The computer readable medium may be a computer readable signal medium or a computer readable storage medium. A computer readable storage medium may be, for example, but not limited to, an electronic, magnetic, optical, electromagnetic, infrared, or semiconductor system, apparatus, or device, or any suitable combination of the foregoing. More specific examples (a non-exhaustive list) of the computer readable storage medium would include the following: an electrical connection having one or more wires, a portable computer diskette, a hard disk, a random access memory (RAM), a read-only memory (ROM), an erasable programmable read-only memory (EPROM or Flash memory), an optical fiber, a portable compact disc read-only memory (CD-ROM), an optical storage device, a magnetic storage device, or any suitable combination of the foregoing. In the context of this document, a computer readable storage medium may be any tangible medium that can contain, or store a program for use by or in connection with an instruction execution system, apparatus or device.\nThe flowchart and block diagrams in the Figures illustrate the architecture, functionality and operation of possible implementations of systems, methods and computer program products according to various embodiments. In this regard, each block in the flowchart or block diagrams may represent a module, segment or portion of code, which comprises one or more executable instructions for implementing the specified logical function(s). It should also be noted that, in some alternative implementations, the functions noted in the block may occur out of the order noted in the figures. For example, two blocks shown in succession may, in fact, be executed substantially concurrently, or the blocks may sometimes be executed in the reverse order, depending upon the functionality involved. It will also be noted that each block of the block diagrams and/or flowchart illustration, and combinations of blocks in the block diagrams and/or flowchart illustration, can be implemented by special purpose hardware-based systems that perform the specified functions or acts, or combinations of special purpose hardware and computer instructions.\nIn view of the foregoing, the scope of the present disclosure is determined by the claims that follow."}
{"text": "The present invention relates to the field of aircraft propeller blades, such as those present on turboprops.\nPropeller blades for turboprops are generally made of metal material. Although propeller blades made of metal material present good strength, they nevertheless present the drawback of being relatively heavy.\nIn order to obtain lighter blades, it is known to use propeller blades that are made of composite material, i.e. by making structural parts out of fiber reinforcement and a resin matrix.\nDocuments US 2013/0017093 and WO 2012/001279 describe making a propeller blade from a streamlined fiber structure into which a portion of a spar is inserted for the purpose of shaping the airfoil portion of the blade so as to form a propeller preform that is subsequently densified with a matrix. The fiber structure, which is made as a single piece by three-dimensional weaving, includes a non-interlinked zone serving to form a housing inside the fiber structure into which a shaping portion of the spar is inserted.\nThe propeller blade as obtained in that way presents both lower overall weight and considerable mechanical strength as a result of having a skin made of a composite material structure (i.e. fiber reinforcement densified by a matrix).\nIn order to impart good strength to the most exposed portion of the blade, i.e. its leading edge, the fiber structure does not include any opening or non-interlinked zone in those portions of the fiber structure that are to form the leading edge of the blade. The non-interlinked zone provided inside the fiber structure opens out both into the bottom edge and into the rear edge of the structure so as to form an opening that enables the shaping portion of the spar to be inserted into the inside of the fiber structure.\nNevertheless, in certain circumstances, such as for example when the blade is subjected to high levels of mechanical loading, to impacts, or to shocks, keeping the shaping portion in its reference position inside the densified fiber structure can be difficult, in particular at the rear edge of the fiber structure which is of reduced strength because of the presence of non-interlinking at this location in the fiber structure."}
{"text": "The present invention relates to high temperature superconductors and methods of making same.\nRecently, there has been much work done on the utilization of ternary oxides containing rare earth elements. barium, and copper for superconductivity above the temperature of liquid nitrogen. To date, these systems typically require four elements: a rare earth; an alkaline earth; copper; and oxygen. Superconductor systems based on rare earth elements have not been entirely satisfactory. The use of rare earth elements in the production of superconductors increases the cost of the resultant superconductors because these rare earth elements are in relatively short supply and are expensive. Furthermore, these systems typically exhibit limited transition temperatures at or below 93 K. Partial substitutions of elements for the member(s) in these systems have not produced significantly higher temperature superconductors.\nTypical high temperature superconductive systems based on rare earth elements also suffer the disadvantage of not being producible in an expedient manner. Therefore, typical methods for producing these high temperature superconductive systems do not provide low cost processing and manufacturability.\nAccordingly, there is a need for a new superconductive system and process for making same."}
{"text": "The efficiency with which a computing device can perform a task used to be, in large part, determined by the frequency with which its Central Processing Unit (CPU) could execute computer-readable instructions. CPUs that executed instructions at a higher frequency could perform a greater amount of work within a given time. As the frequency with which CPUs executed instructions began to reach practical limits, computing devices began to include multiple CPUs that could perform multiple tasks in parallel. For such computing devices, if a task was designed, and programmed, such that segments of computer-executable instructions could be executed in parallel, then that task could be performed more efficiently by executing, in parallel, those segments on the multiple CPUs of the computing device.\nUnfortunately, for many tasks, the input to one set of computer-executable instructions is based on, or affected by, the output of a prior set of computer-executable instructions. Consequently, such tasks can be difficult to execute in parallel and the execution of such tasks may utilize only one of the multiple CPUs that can be available. Nevertheless, a computing device with multiple CPUs may still provide efficiency benefits, even for tasks that are difficult to parallelize, because users often request the execution of multiple such tasks in parallel, thereby parallelizing the achievement of the user's overall goals.\nTo provide purchasers of computing devices with expandability, add-on hardware that itself contains one or more CPUs can be added to the computing device. Similarly, virtual computing environments can be used and expandability can, thus, be achieved via the addition of further virtual computing environments. In such cases, however, it may not be practical to extend core software, such as the operating system, to the multiple CPUs or other execution hardware. Instead, the each of the execution environments that were subsequently added can execute their own copy of an operating system or other core software. While such an operating system may be identical to the operating system executing on the original “host” CPU or execution environment, communication between the two, or more, copies of the operating system will be limited to, essentially, inter-computing-device communications. Thus, the addition of such additional computing environment, virtual or physical, to a computing device, can expand the ability of the computing device to execute multiple independent tasks in parallel, but may not expand the ability of the computing device to execute a single, parallelizable, task in a more efficient manner."}
{"text": "1. Field of the Invention\nThe invention relates to depth filter media.\nIn this specification \"depth filter\" means a filter with pores capable of removing from a fluid particles that may be smaller than the size of pores in the filter, the particles being trapped by progressive interception during changes of direction of the pores. Such depth filters have a high dirt capacity.\nDepth filters are commonly utilized in the form of a thick continuous cylinder of filter medium that surrounds a central core and may be provided with an external cage. Such depth filters have, because of their low external area and great depth (typically about 15 mm) a relatively high pressure drop when a fluid is passed through the filter. On the other hand, as mentioned above, such filters have relatively high dirt capacity, because of the high internal void space in which contaminant can accumulate.\n2. Brief Review of the Prior Art\nGB-A-585295 discloses a filter element formed by an elongated web of cellulose filter material folded to extend back and forth between the inner and outer surfaces of the tubular body. The folds extend longitudinally and a tubular binding strip is secured to the inner folds. No drainage layer is provided. GB-A-1389199 discloses a filter formed by corrugating a flat sheet felt medium backed by woven wire. The ends of the sheet are sealed after pleating. The woven wire supports the felt and no drainage layer is provided.\nGB-A-1460925 discloses a filter element formed from a corrugated sheet of depth filter medium in which the edges of the sheet are connected together after corrugation. Separate upstream and downstream protective layers are provided. U.S. Pat. No. 4,233,042 discloses a pleated filter medium formed from sheets of glass fibre. EP-A-0083789 discloses a pleated sheet of filter material such as a fibrous filter material. The edges cf the sheet are side-sealed after corrugation.\nAlthough pleating increases the surface area per unit volume of filter material, and thus increases the dirt capacity, the presence of side seals formed after pleating can be disadvantageous. First, the presence of the seal forms an obstruction to uniform flow through the filter element. Secondly, the formation of the side seal requires additional manufacturing steps after pleating (see GB-A-1460925). Thirdly, the presence of a side seal forms a potential leakage path because complete sealing is difficult.\nIn order to achieve optimum performance of a pleated filter, it is necessary to provide a relatively coarse upstream drainage layer to allow drainage of fluid down between the pleats and to give a void space for the accumulation of dirt. They also support the filter medium. It is also necessary to provide a relatively coarse downstream drainage layer to allow filtrate to drain between the pleats from an inner filtration layer and also to support the filter medium against applied pressure.\nThe drainage layers are conventionally formed separately from the filter medium by layers of non-woven fabrics or nets located on either side of the filter medium (see GB-A-1460925). The provision of such separate layers and their attachment to the filter medium complicates the manufacture of the filter and increases its cost."}
{"text": "As the size of transistors is reduced to nanometer scale dimensions, for example in the form of ultra-thin body (UTB) silicon-on-insulator (SOI) field effect transistors (FETs), FinFETs and nanowire FETs, the unwanted resistance associated with transistor sources and drains becomes an ever increasing burden on the performance of these devices and of the integrated circuit products manufactured using such transistors. Furthermore, a reduction of dopant activation is predicted theoretically and demonstrated experimentally when the transistor source and drain regions are reduced in size below approximately 10 nm. By dopant activation, we mean desired free carrier (electron or hole) contributions from deliberately introduced impurity species in a semiconductor host. This decrease in nanoscale dopant activation further contributes to undesirably high resistance of doped source/drain (S/D) regions both at the nanoscale metal contacts and within the bulk portion of the nanoscale doped regions. The resistance of metal contacts to a semiconductor increases if effective doping in the semiconductor decreases, the increase being primarily due to the presence of a Schottky barrier at metal-semiconductor contacts.\nIt is known that a high concentration of doping in a shallow region of semiconductor proximate to a metal-semiconductor interface can reduce the resistance of the metal-semiconductor contact by decreasing the width of the Schottky barrier. Although it is the barrier width that is reduced, from an electrical response point of view (for example current—voltage measurement), it appears that the Schottky barrier height is reduced. An early article that describes this “effective barrier height” reduction by surface doping is by J. M. Shannon, “Control of Schottky barrier height using highly doped surface layers” in Solid-State Electronics, Vol. 19, pp. 537-543 (1976). It is also known that a high concentration of dopant atoms can be introduced into a shallow region of a semiconductor proximate to a metal contact by so-called dopant segregation out of a metal silicide. A. Kikuchi and S. Sugaki reported in J. Appl. Phys., Vol. 53, No. 5, (May 1982) that implanted phosphorus atoms piled up near a PtSi—Si interface during PtSi formation and reduced the measured height of the Schottky barrier to n-type silicon. The reduction of the measured (effective) barrier height of the Schottky diode was attributed to piled up phosphorus atoms in the silicon causing the barrier to be more abrupt. That is, the result was attributed to the effect described by Shannon in 1976.\nFor the past several decades the silicon microelectronics industry has relied on high doping concentrations in the silicon proximate metal-silicon contacts as a means of obtaining acceptably low contact resistances to transistor sources and drains. The contact metal has for the most part been a metal silicide, most recently nickel silicide or nickel platinum silicide. This approach to minimizing contact resistance is expected to be insufficient in the future as the transistor dimensions continue to shrink and the contact resistance becomes a larger portion of the total resistance between the source and drain (hence becoming a serious performance-limiting factor). The most recent International Technology Roadmap for Semiconductors (ITRS), published in 2011, reports that there is no known solution to the contact resistance problem in bulk MOS transistors when the transistor gate length scales to 18 nm, as expected in year 2014, and a specific contact resistance of no more than 1.0×10−8 Ohm·cm2 is specified. It is becoming increasingly apparent that the Schottky barrier at metal-semiconductor contacts must be reduced in order to reduce the contact resistance to acceptable levels, i.e. well below 1.0×10−8 Ohm·cm2 in the case of MOS transistor doped source/drain contacts. A technology that is capable of reducing the Schottky barrier and hence reducing the resistance of contacts to doped semiconductor regions may also be applied to so-called “metal source/drain transistors” which do not have doped source and drain but rather utilize a direct contact between the metal and the transistor channel (the region of free carriers that are modulated by the electrical potential on a gate and that transport current between the source and the drain).\nA body of work published in 1991-1992 reported experimental verification of theoretical predictions by Baroni, Resta, Baldereschi and others that a double intralayer formed by two different elements would create an interface dipole, capable not only of modifying heterojunction band discontinuities, but also of generating band discontinuities in homojunctions. McKinley et al. first reported obtaining 035-0.45 eV band offsets at {111}-oriented Ge homojunctions using Ga—As dipole intralayers in a 1991 article “Control of Ge homojunction band offsets via ultrathin Ga—As dipole layers”, J. Vac. Sci. Technol. A 9 (3), May/June 1991 and in a similar article in 1992 “Control of Ge homojunction band offsets via ultrathin Ga—As dipole layers”, Applied Surface Science Vol. 56-58, pp. 762-765 (1992).\nArsenic, gallium, and germanium depositions were done at room temperature on p-type Ge(111) substrates. Valence band offsets were measured by in situ core level x-ray photoluminescence. The deposited Ge region (overlayer) had a valence band offset to the Ge substrate as manifested by a splitting of the Ge 3d core level into two components; one due to the Ge substrate and the other to the Ge overlayer. Both positive and negative valence band offsets were obtained in Ge homojunctions by introducing Ga—As dipole intralayers with either “Ga-first” or “As-first” growth sequences. The band offset was found to be 0.35-0.45 eV with the Ge valence band edge on the As side of the junction at a lower energy (i.e., more bound). Dipole intralayers were explained on the basis of the Harrison “theoretical alchemy” model described by W. A. Harrison et al. in “Polar Heterojunction Interfaces”, Phys. Rev. B 18, 4402 (1978). Intralayer control of band discontinuities was thus applied to homojunctions, expanding the potential domain of band offset engineering beyond semiconductor heterojunctions.\nIn 1992, Marsi et al. followed up on the reports by McKinley et al. with the articles “Microscopic manipulation of homojunction band lineups”, J. Appl. Phys., Vol. 71, No. 4, 15 Feb. 1992, “Homojunction band discontinuities induced by dipolar intralayers: Al—As in Ge”, J. Vac. Sci. Technol. A 10(4), July/August 1992 and “Local nature of artificial homojunction band discontinuities”, J. Appl. Phys. 72 (4), 15 Aug. 1992. In the first article, Marsi et al. reported valence-band discontinuities at Si—Si and Ge—Ge homojunctions when III-V double intralayers of atomic thickness were inserted at the interfaces. Valence band discontinuities were again measured by in situ core level x-ray photoluminescence. In Ge samples, a deposited Ge region (overlayer) had a valence band offset to a Ge substrate as evidenced by a splitting of the Ge 3d core level into two components and a deposited Si region had a valence band offset to a Si substrate as evidenced by a splitting of the Si 2p core level. The observed discontinuities with magnitudes in the range 0.4 to 0.5 eV (for example 0.5 eV for Si—P—Ga—Si and 0.4 eV for Si—P—Al—Si) were in qualitative agreement with theoretical predictions although most theories estimate larger valence band discontinuities due to the dipole effect. A III-V intralayer at a group-IV homojunction systematically induced an artificial valence-band discontinuity when the anion was deposited first. It was also reported that in the case of Si—Si homojunctions with Al—P or Ga—P intralayers, a reversal of the interface deposition sequence led to a reversal of the valence-band discontinuity, as expected.\nIn the second article it was shown, again using x-ray photoemission, that a similar band offset effect can be induced using Al—As as a “dipolar intralayer” between two regions of {111}-oriented germanium. Specifically, an offset of 0.4 eV was obtained for the “anion-first” Ge(substrate)-As—Al—Ge(overlayer) sequence, consistent with the “anion-first” As—Ga sequence reported by McKinley, the overlayer component exhibiting a lower binding energy with respect to the substrate component. In the third article, multiple III-V bilayer (intralayer) stacks were investigated. The measured value of valence band offset remained the same, 0.5 eV, for an individual double layer, for double-stacked bilayers and for triple-stacked bilayers. Experiments performed on 2(Ga—P) and 2(P—Ga) were fully consistent with those on 2(Al—P) and 2(P—Al); no substantial increase was observed on going from the individual bilayers to two bilayers or even to three bilayers. It was therefore concluded that stacked interfacial III-V bilayers do not increase the effect of an individual bilayer, contrary to elementary predictions based on sequential dipoles.\nIn U.S. Pat. Nos. 7,084,423, 7,176,483, 7,462,860, and 7,884,003 and in pending U.S. patent application 2011/0169124, Grupp and Connelly described metal-semiconductor contacts having an interfacial layer at the interface between a metal and a group IV semiconductor for the purpose of reducing the Schottky barrier at the contact and, hence, reducing the specific resistivity of the contact. A monolayer of arsenic (or nitrogen) was included amongst the possible embodiments/specifications of the interfacial layer."}
{"text": "Boron neutron capture therapy (BNCT) is a unique tumor cell targeting therapy (Soloway et al, Chem. Rev. 98:1515-1562, 1998; Barth et al., Clinical Res. 11:3987-4002, 2005). It is one of the binary cancer treatment systems that are based on the selective accumulation of boron-10 (10B) in tumors followed by irradiation with a neutron source. The selective accumulation of 10B in tumors and the subsequent capture of an epithermal neutron by a 10B atom, which produces an α-particle (He2+) and a lithium nucleus (7Li+) ejected in opposite directions within the tumor mass, make BNCT an attractive therapy for targeting tumors. The average track of these densely ionizing particles is approximately 14 μm, about the diameter of one cell, so that the killing of tumor cells is highly efficient (see Barth et al., Cancer 70:2995 3007, 1992; Barth et al., Neurosurg. 44:433 451, 1999; Soloway et al., J. Neurooncol. 33:9-18, 1997; and Soloway et al., 1998, supra). This promising radiation therapy requires the administration of tumor-seeking compounds containing 10B atoms that concentrate preferentially in tumor cells prior to irradiation with suitable low energy neutron beams. This binary modality is a highly attractive cancer therapy in that neither the thermal neutrons nor the boron carrier molecule has significant cytotoxic effect alone, but in combination they produce highly radiobiologically effective particles. The alpha particle and Li ion that are released as fission products when a thermal neutron is absorbed by a 10B atom have high linear energy transfer (LET) and a range similar to the dimensions of a mammalian cell. Thus, the absorbed dose, which is potent, is confined to cells adjacent to the boron atoms. The potential exists, therefore, to destroy tumor cells by radiosurgery at the microscopic level while sparing normal tissue in the vicinity of the tumor. BNCT is especially attractive for malignant brain tumors because it targets and destroys malignant cells but spares normal cells, thus preventing undesirable side effects common in standard chemo- or radiotherapy.\nIn recent years, several research groups have developed a variety of 10B carriers, including porphyrin derivatives with improved tumor selectivity over prior boron neutron capture agents, such as disodium mercapto-closo-dodecaborate (BSH) and L-4-dihydroxy-borylphenylalanine (BPA). Porphyrin derivatives are currently being tested in clinical trials in the U.S., Europe, and Japan for the treatment of patients with glioblastomas and melanomas (see, e.g., Bonnett et al., Chem. Soc. Rev. 24:19-33, 1995; Kageji et al., J. Neurooncol. 33:117-130, 1997; Pignol et al., Br. J. Radiol. 71:320-323, 1998; Elowitz et al., Neurosurgery 42:463-469, 1998). Although BSH and BPA have been shown to be safe and efficacious in animal models, several problems remain. For instance, BSH is sensitive to air-oxidation (Tolpin et al., Inorg. Chem. 17:2867-2873, 1978) and both BSH and BPA have low retention times in tissues and only moderate selectivity for tumor cells (Capala et al., Radiation Res. 146:554-560, 1996). At present, clinical progress is stalled due to the absence of effective boron delivery agents that target tumor cells with high selectivity. Obviously, the ultimate success of BNCT will be dependent upon whether adequate concentrations of boron neutron capture agents and low-energy neutrons can be selectively and effectively delivered to tumor cells.\nHypoxic tumor cells that comprise up to 50% of the mass of viable cells in human tumors pose a special impediment to effective radio- and chemotherapy. In particular, hypoxic cells are highly radio-resistant and capable of proliferation following radiation treatment. If boron containing therapeutics that interact with epithermal neutrons are not delivered to hypoxic sites in tumors it can be anticipated that the effectiveness of BCNT will be significantly compromised. In principal, BNCT is well suited to treating poorly oxygenated tumors since reduced radio-sensitivity in the absence of oxygen is less pronounced for high LET radiations, yet the effectiveness of BNCT therapy is dependent on delivery of the agent carrying the 10B atom to the tumor site.\n2-nitroimidazoles readily penetrate tumors and can reach intratumor concentrations approaching 1 mM. As Varghese et al. showed in 1976 (Scobie et al., 1994), 2-nitroimidazoles also undergo nitroreductive activation under hypoxic conditions to yield electrophilic species that form adducts with cellular macromolecules, such as DNA and proteins. These adducts are retained in hypoxic cells and represent a useful targeting mechanism. In fact, carboranes have been linked to nitroimidazoles (Wood et al., 1996; Swenson et al., 1996) in attempts to target boron selectively to the hypoxic regions of tumors, but these boron delivery agents suffered greatly from poor water solubility. Thus, there remains a need for highly effective reagents for use in treating tumor tissue."}
{"text": "1. Field of the Invention\nThis invention relates to a simplified, potential low cost system for drywall tape installation.\nSince the tape is fully imbedded after only one operation, there is a reduction in labor to complete the drywall taping procedure.\n2. Background Art\nThe repetitive and laborious process of taping drywall has lead to several systems of taping, each of which have serous undersirable limitations.\nSince the advent of drywall technology, the taping of the seams created when sheets of drywall are nailed in place has been a pedestrian exercise in handling the paper tape and adhesive wet mud. The components are currently sold separately and are combined only upon actual application with all prior art.\nIn the time consuming manual application of the taping materials, the adhesive mud is manually scooped out of a container and applied with a smearing action to drywall seams with a drywall knife. Then dry paper tape is un-rolled and laid upon the adhesive surface and wiped down onto the wet mud. The tape is later top coated with more mud and wiped smooth after each coat to flatten the surface. The materials are allowed to dry between coats.\nOne layer of mud goes under the tape to achieve adhesion and at least two layers of mud are commonly applied over the tape to \"embed\" the tape within the mud. The mud is allowed to dry between coats. Mud is applied to the drywall seams and wiped down until the mud appears smooth and the gap between the drywall panels are made flat.\nThe predominate commercial taping system in use, is the Ames system. The Ames system consist of a series of tools designed to handle both dry, non-adhesive paper tape and wet adhesive mud to make the application procedure faster than manual installation.\nThe Ames system tape tools consist of; a manual mud pump, a taper tool, roller, finishers and a standard drywall knife. Many of the Ames tools are complex, rather expensive and require extensive training, adjustment and manual dexterity to operate.\nIn the Ames system the drywall mud is pumped out of a bucket into a cylinder within the Ames taper. The Ames taper tool holds a roll of dry paper tape. When the cylinder is full of mud the taper is held against the seam and the taper head is rolled along the seam. The cylinder is emptied of the mud it holds, as a wheel on the head which contains a pulley, reels up a steel cable attached to a piston within the cylinder. The mud then exits the tool out a hole at the head as dry paper tape is rolled onto the mud as it is ejected upon the wall.\nThe seam must be \"wiped down with a drywall knife to complete the taping operation with the Ames system. The seam then requires a top coat of mud to fully imbed the tape.\nWith the Ames system full embedding occurs only after three operations; Tape and mud application to the drywall, wiping the tape down and then top coating the seam.\nThe previously mentioned tools have various limitations and short comings. The Ames system is very expensive as it requires many unique part be machined and the tools maintained in proper adjustment. The taper tool is heavy and hard to use. The Ames taper requires that many precision parts operate in the hostile environment of gritty, adhesive viscous drywall mud. The tool must be cleaned, oiled and adjusted often to not breakdown or become jammed and stop working. Parts wear out and break requiring substantial maintenance and expense. The Ames taper does not apply mud to the top side of the drywall tape.\nA more primitive taping tool called the \"banjo\" utilizes the concept of a reel of dry tape that is threaded though a chamber of wet mud as it exits the tool and the muddy tape is then directly applied to the drywall for manual wipe down. As the banjo concentrates all the mud under the tape, very little mud ends up on the back side of the tape. With a banjo taper another coat of mud must be applied to fully embed the tape.\nThe banjo is slow and inefficient as the mud may not properly and fully coat the underside of the tape, leaving a dry blister under the tape. The banjo is heavy and awkward to use and requires that the muddy tape be wiped down after it is unrolled upon the drywall seam. The banjo will not readily do inside corners and does not apply mud to both sides of the drywall tape.\nIt would be desirable to have an effective labor saving system that is inexpensive to manufacture, easy to use and light weight. The system should be able to operate in the hostile environment of the materials and not be adversely affected by the mud drying and or a failure to thoroughly clean the tools after each use.\nThe system should utilize current materials and procedures common to the drywall taping industry. Use of the tools should be natural and easy to learn. It would be desirable to reduce the occurrence of blisters of air which occasionally occur under the tape. The system should fully embed the drywall tape in the initial application; thus reducing the number of required steps and eliminating some labor."}
{"text": "The invention relates to a method for producing a semiconductor structure having a plurality of vertical and at least one lateral semiconductor component integrated in a disk-shaped semiconductor body having a first surface and a second surface opposite the first surface, with the lateral semiconductor components being disposed beneath the first surface, and the plurality vertical semiconductor components extending between the first and the second surfaces. The at least one lateral semiconductor component being disposed inside a partial structure which is insulated from adjacent regions. The thickness of the semiconductor body is reduced in the region of the partial structure. As used herein, the term semiconductor component is intended to mean semiconductor device or devices integrated into a semiconductor body.\nA method of producing, such a semiconductor structure is disclosed in EP-A-0,193,172. In this instance, a vertical MOS transistor is described which is disposed in a semiconductor body along with a lateral switching circuit. The MOS transistor and the lateral switching circuit are isolated from one another by a p-n junction. It does not follow from EP-A-0,193,172 to integrate a plurality of vertical components in the semiconductor body. A disadvantage in the disclosed type of insulation EP-A-0,193,172 is that the space-charge region present at the p-n junction requires a great deal of space that cannot be used for the insertion of components. Moreover, the doping of the regions lying adjacent to the isolation cannot be freely selected; namely--for obtaining a space-charge region--a p-doped and an n-doped zone must always be present alternatingly so that a p-n junction results. Furthermore, it is disadvantageous that, at a p-n junction, the correct polarity (in the blocking direction) must be heeded, because otherwise the isolating effect of the junction is lost. The permitted potential in the one region is therefore a function of that in the other region--located behind the p-n junction."}
{"text": "The invention relates to a method for the manufacture of a wire electrode for the spark-erosion process and a wire electrode for this method, which consists of a core of steel, an intermediate layer of copper or a high copper-containing alloy, and an outer layer with at least 40% zinc or a core of steel and a zinc alloy outer layer having a zinc content of from 40-60%.\nWire electrodes, which are designed with multiple layers, which have a core of steel, an intermediate layer of copper arranged around the core, and an outer zinc-containing layer, are known, for example, from the DE-PS 29 06 245. All known wire electrodes, which are constructed with a steel core, do indeed have an increased strength compared with copper or brass electrodes, however, all of these erosion electrodes have the disadvantage that if they have comparatively high strengths, their electrical conductivity is very low and merely reaches 8 Sxc2x7m/mm2. This comparatively high tensile strength is particularly advantageous for the cutting of high or thick workpieces or, however, also for the cutting of very small parts since high wire tensions are here demanded. However, these known wire electrodes have the disadvantage that their erosion performance is relatively low. In particular, in the case of cutting very small parts with erosion-wire diameters of up to 10 xcexcm, high wire strengths are demanded in order to reduce deflection and vibration of the wire to a minimum. Tungsten or molybdenum wires have been used up to now for the cutting of very small parts with wire thicknesses of 100 xcexcm and less, however, tungsten or molybdenum wires are extremely expensive to manufacture. Erosion wires with a steel core and a brass outer layer have been unable to be successful up to now for this purpose since they always, in comparison to tungsten wires, showed a lesser strength and a poorer erosion behavior.\nThe basic purpose of the invention is to provide high-strength erosion electrodes with a core of steel of the abovementioned type and a method for their manufacture, which electrodes have strengths which are 1800 N/mm2 and higher and, in addition, have a comparatively high conductivity of 10 Sxc2x7m/mm2 and more."}
{"text": "For example, an Li ion battery is known as a high-voltage and high-energy density battery. The Li ion battery is a cation-based battery utilizing a reaction between an Li ion and a cathode active material and a reaction between an Li ion and an anode active material. Meanwhile, as an anion-based battery, a fluoride ion battery utilizing a reaction of a fluoride ion is known.\nAlso, it is known that an ionic liquid is used for a liquid electrolyte. For example, in Patent Literature 1, a nonaqueous electrolyte secondary battery provided with an electrolyte having an ionic liquid and a lithium salt is disclosed. In addition, as the ionic liquid, an ionic liquid containing fluoride complex anions such as 1-butylpyridinium tetrafluoroborate is disclosed.\nAlso, in Patent Literature 2, a fluoride ion battery provided with an electrolyte dissolving a fluoride salt in a solvent and an addition agent including a fluoride composite producing species is disclosed. In addition, examples of the solvent include an ionic liquid (claim 14). Also, examples of the fluoride composite producing species include tris(pentafluorophenyl)borane (claim 3)."}
{"text": "1. Field of the Invention\nThe present invention relates to a lamp and, more particularly, to a lamp such as a wall lamp, night lamp and the like.\n2. Description of the Related Art\nA conventional lamp comprises a base, a lamp socket mounted on the base, a light emitting member mounted on the lamp socket, and an outer shade mounted on the base to surround the light emitting member. Thus, the base is mounted on a wall (or ceiling) to attach the lamp onto the wall to provide a lighting effect. The outer shade of the lamp is provided with patterns to enhance the outer appearance of the lamp. However, the outer shade of the lamp has a fixed structure so that the patterns of the outer shade lack variation, thereby limiting the versatility of the lamp."}
{"text": "Traditional methods for detecting biological compounds in vivo and in vitro rely on the use of radioactive markers. For example, these methods commonly use radiolabeled probes such as nucleic acids labeled with 32P or 35S and proteins labeled with 35S or 125I to detect biological molecules. These labels are effective because of the high degree of sensitivity for the detection of radioactivity. However, many basic difficulties exist with the use of radioisotopes. Such problems include the need for specially trained personnel, general safety issues when working with radioactivity, inherently short half-lives with many commonly used isotopes, and disposal problems due to full landfills and governmental regulations. As a result, current efforts have shifted to utilizing non-radioactive methods of detecting biological compounds. These methods often consist of the use of fluorescent molecules as tags (e.g. fluorescein, ethidium, methyl coumarin, rhodamine, and Texas red), or the use of chemiluminescence as a method of detection. Presently however, problems still exist when using these fluorescent and chemiluminescent markers. These problems include photobleaching, spectral separation, low fluorescence intensity, short half-lives, broad spectral linewidths, and non-gaussian asymmetric emission spectra having long tails.\nFluorescence is the emission of light resulting from the absorption of radiation at one wavelength (excitation) followed by nearly immediate reradiation usually at a different wavelength (emission). Fluorescent dyes are frequently used as tags in biological systems. For example, compounds such as ethidium bromide, propidium iodide, Hoechst dyes, and DAPI (4′,6-diamindino-2-phenylindole) interact with DNA and fluoresce to visualize DNA. Other biological components can be visualized by fluorescence using techniques such as immunofluorescence which utilizes antibodies labeled with a fluorescent tag and directed at a particular cellular target. For example, monoclonal or polyclonal antibodies tagged with fluorescein or rhodamine can be directed to a desired cellular target and observed by fluorescence microscopy. An alternate method uses secondary antibodies that are tagged with a fluorescent marker and directed to the primary antibodies to visualize the target.\nAnother application of fluorescent markers to detect biological compounds is fluorescence in situ hybridization (FISH). This method involves the fluorescent tagging of an oligonucleotide probe to detect a specific complementary DNA or RNA sequence. An alternative approach is to use an oligonucleotide probe conjugated with an antigen such as biotin or digoxygenin and a fluorescently tagged antibody directed toward that antigen to visualize the hybridization of the probe to its DNA target. FISH is a powerful tool for the chromosomal localization of genes whose sequences are partially or fully known. Other applications of FISH include in situ localization of mRNA in tissues samples and localization of non-genetic DNA sequences such as telomeres.\nFluorescent dyes also have applications in non-cellular biological systems. For example, the advent of fluorescently-labeled nucleotides has facilitated the development of new methods of high-throughput DNA sequencing and DNA fragment analysis (ABI system; Perkin-Elmer, Norwalk, Conn.). DNA sequencing reactions that once occupied four lanes on DNA sequencing gels can now be analyzed simultaneously in one lane. Briefly, four reactions are performed to determine the positions of the four nucleotide bases in a DNA sequence. The DNA products of the four reactions are resolved by size using polyacrylamide gel electrophoresis. With singly radiolabeled (32P or 35S) DNA, each reaction is loaded into an individual lane. The resolved products result in a pattern of bands that indicate the identity of a base at each nucleotide position. This pattern across four lanes can be read like a simple code corresponding to the nucleotide base sequence of the DNA template. With fluorescent dideoxynucleotides, samples containing all four reactions can be loaded into a single lane. Resolution of the products is possible because each sample is marked with a different colored fluorescent dideoxynucleotide. For example, the adenine sequencing reaction can be marked with a green fluorescent tag and the other three reactions marked with different fluorescent colors. When all four reactions are analyzed in one lane on a DNA sequencing gel, the result is a ladder of bands consisting of four different colors. Each fluorescent color corresponds to the identity of a nucleotide base and can be easily analyzed by automated systems.\nThere are chemical and physical limitations to the use of organic fluorescent dyes. One of these limitations is the variation of excitation wavelengths of different colored dyes. As a result, simultaneously using two or more fluorescent tags with different excitation wavelengths requires multiple excitation light sources. This requirement thus adds to the cost and complexity of methods utilizing multiple fluorescent dyes.\nAnother drawback when using organic dyes is the deterioration of fluorescence intensity upon prolonged exposure to excitation light. This fading is called photobleaching and is dependent on the intensity of the excitation light and the duration of the illumination. In addition, conversion of the dye into a nonfluorescent species is irreversible. Furthermore, the degradation products of dyes are organic compounds which may interfere with biological processes being examined.\nAnother drawback of organic dyes is the spectral overlap that exists from one dye to another. This is due in part to the relatively wide emission spectra of organic dyes and the overlap of the spectra near the tailing region. Few low molecular weight dyes have a combination of a large Stokes shift, which is defined as the separation of the absorption and emission maxima, and high fluorescence output. In addition, low molecular weight dyes may be impractical for some applications because they do not provide a bright enough fluorescent signal. The ideal fluorescent label should fulfill many requirements. Among the desired qualities are the following: (i) high fluorescent intensity (for detection in small quantities), (ii) a separation of at least 50 nm between the absorption and fluorescing frequencies, (iii) solubility in water, (iv) ability to be readily linked to other molecules, (v) stability towards harsh conditions and high temperatures, (vi) a symmetric, nearly gaussian emission lineshape for easy deconvolution of multiple colors, and (vii) compatibility with automated analysis. At present, none of the conventional fluorescent labels satisfies all these requirements. Furthermore, the differences in the chemical properties of standard organic fluorescent dyes make multiple, parallel assays quite impractical since different chemical reactions may be involved for each dye used in the variety of applications of fluorescent labels."}
{"text": "1. Field of the Invention\nThe present invention relates to a structure of ESD protection circuits on a BEOL layer and a method of fabricating the same.\n2. Description of the Prior Art\nThere are various types of circuits which protect against Electrostatic Discharge (ESD). ESD means a discharge from an electrostatically charged person or machine to a semiconductor device, a discharge from a charged semiconductor device to earth, and the like. When ESD occurs in a semiconductor device, a large amount of electric charge flows from the terminals of the semiconductor device into the semiconductor device, which raises the voltage inside the semiconductor device and may cause a dielectric breakdown of the internal circuits of the semiconductor device.\nConventional ESD protection circuits are disposed at a same level of the internal circuits. These circuits occupy too much space, however, which influences the potential miniaturization of the semiconductor device."}
{"text": "In the field of electronic devices, such as systems on a chip, a phase lock loop is often used to generate a master clock signal. That master clock signal is then divided downward into a variety of different clock signals for use by different components within the system on chip.\nThis division is performed by coupling multiple divider circuits to the output of the phase locked loop. For example, a divide by 2, a divide by 4, and a divide by 8 divider circuit are each coupled to the phase locked loop. However, this approach has drawbacks. For example, power consumption may be undesirably high.\nIt is therefore desirable for new divider circuits to be developed."}
{"text": "A great deal of time, money, and effort is spent in the fight against security risks, such as viruses, malware, exploits, intrusions, and other dangers to enterprise computing. A malware infection may require a compromised system to be erased and re-imaged to remove the malware, which can also result in great expenditures of time and money as well as loss of productivity.\nWhile terms such as viruses, malware, and exploits may have specific definitions in the art, as used herein the term “malware” shall refer to any type of computer code and/or other data introduced from outside a computer system or network of computers that may be used to damage or otherwise cause aspects of the computer system or network of computers to malfunction or perform in an unintended manner. Thus, as used herein, the term “malware” encompasses viruses, exploits, worms, and any type of programmatic security risk.\nMalware may be introduced into a computer network through a wide variety of avenues, such as email or a malformed JPEG or other code on a web page that is downloaded to a computer via a web browser. Malware generally follows a common lifecycle; namely, the malware will exploit a weakness in the computer system or network, the malware will drop and execute a payload, the malware will attempt to escalate its privileges, the malware will take steps to persist on the system and/or network despite efforts to remove it, and finally the malware will propagate to other computers. If the malware is unable to be detected using current approaches, then countermeasures cannot be enacted and systems may become infected."}
{"text": "1. Field of the Invention\nThe invention relates to a time switch whereby the on and off states of internal contact points control the operation of an electrical load according to a time schedule.\n2. Description of the Background Information\nA conventional time switch, such as the type described in Japanese Kokai (laid open) Patent S63-29427, includes a modular structure assembled from a motor that rotates at a uniform speed, a ratchet gear driven by a gearwheel transmission powered by the motor, and a first gear that meshes with the ratchet gear. The time switch also includes a sub-dial that rotates clockwise at a rate of one revolution per hour relative to the rotation of the ratchet gear, a gear that meshes with a second gear attached to the rotating shaft of the sub-dial, a dial that rotates clockwise once during a 24-hour period relative to the rotation of the sub-dial, a scheduling clip with a finger that projects from the circumference of a dial used to set a desired operating time schedule, and an internal contact point activating mechanism that opens and closes the contact points through contact with a contact element that is freely settable to a desired time. In the time switch of Japanese Kokai (laid open) Patent S63-29427, the ratchet gear, which is rotatably driven by motor torque supplied through the gearwheel transmission, rotatably drives the sub-dial through the sub-dial shaft. Although the dial is rotationally driven by meshing with the second gear on the sub-dial shaft, the sub-dial shaft rotates only in one direction due to the ratchet mechanism formed by the first gear and ratchet gear at the rotating sub-dial shaft. As a result, even though the sub-dial can be turned manually to set the time, it can only be turned in one direction, thereby making it time-consuming to set the time. For example, when changing the time schedule due to the enactment of daylight savings time, it would be easy to turn the sub-dial back one hour from a current setting of 7 o'clock to a 6 o'clock setting. Because the sub-dial can only be turned in the clockwise direction, however, it must be advanced 23 hours to effect the desired time change. As a result, changing the time becomes a time-consuming operation."}
{"text": "Steerable satellite antenna assemblies for satellite-to-ground links typically require high gain, low mass, and high reliability. One approach for a steerable satellite antenna assembly is to use a fixed feed source so as to eliminate performance degradations otherwise associated with a moving feed source. Performance degradations for a moving feed source typically include losses due to mechanical rotary joints, RF cable connectors, flexible waveguides, and lengthy cables associated with rotary actuators.\nSteerable satellite antenna assemblies also need to avoid the “keyhole effect,” which is a physical limitation due to the orientation of the antenna rotation axis caused by a limited motion range of the actuators. The keyhole effect causes the antenna to momentarily disrupt communications when reaching its physical limitation so as to allow for the actuators to reposition before resuming steering.\nOne approach for a steerable satellite antenna assembly with a fixed antenna feed is disclosed in U.S. Pat. No. 6,492,955 to Amyotte et al. The antenna assembly includes a subreflector secured to a frame rotatably mounted on a support structure via a first motor, and an antenna feed located at a first focus of the subreflector for illuminating the same. The antenna feed is fixed to the structure and has a feed axis pointing at the subreflector. A parabolic reflector having a focus in common with a second focus of the subreflector to transfer the signal between the same, and a planar reflector is secured to the frame, and has a beam axis. The planar reflector is rotatably mounted on the frame via a second motor to transfer the signal between the parabolic reflector and a target. The antenna assembly includes a controller connected to the motors to steer at the target anywhere within a full spherical angular range.\nAnother approach for a steerable satellite antenna assembly with a fixed antenna feed is disclosed in U.S. Pat. No. 5,198,827 to Seaton. The antenna assembly includes an antenna feed structure for emitting electromagnetic radiation, and a subreflector for redirecting the emitted radiation. A main antenna reflector projects radiation redirected by the subreflector as an antenna beam. A mechanical arrangement rotates the subreflector about the rotation point.\nEven in view of the above steerable satellite antenna assemblies, there is still a need for such assemblies to have a more compact geometry. Continued growth and demand for bandwidth has led to new commercial satellite constellations. For example, the O3b satellite constellation is deployed in a medium earth orbit (MEO) and the OneWeb satellite constellation is to be deployed in a low earth orbit (LEO). A more compact antenna assembly reduces the size and weight of the satellites, as well as costs. In addition, MEO and LEO satellites typically operate at a wider beam scanning range as compared to a traditional geosynchronous orbit. Consequently, another desirable attribute is for the compact antenna assemblies to also operate over a wider scanning range."}
{"text": "This invention relates to pre-lined art board sheet member, and more particularly to an art sheet member which is pre-lined for the preparation of printed advertising media.\nIn the preparation of advertising material for printed advertising, plain, unlined art board is used by the advertising artist to prepare lay-outs, and after approval by the advertiser, the completed advertising illustration and text. The completed advertising material is then submitted to the color separator who prepares film of the different color separations, which in turn is presented to the printer. The printer then places the different colored film into chases or vacuum frames for preparation of the color plates for the printing press.\nThe original borders, center lines, and fold lines are laid out upon the art board for each separate advertising material by the advertising artist. After the printer receives the color films, he must extend the various border center lines and other lines upon the film for registration with his printing plates, in order for the artwork to properly register on the printed pages. A certain amount of time must be spent by the artist for each different advertising lay-out, just for placing the border and center lines upon the art board. Drawing and sketch lines are frequently changed upon the art board, requiring erasure and smudging.\nMoreover, because the artist has to measure and draw the border, center lines and fold lines for each new advertising creation, errors are frequently made which continue into the completed advertising material. Such errors can be quite expensive causing complete re-printing of numerous pages of advertising material, re-assembly of the print chases and printing plates, and even re-preparing the lay-outs and completed advertising drawings. Such errors can create considerable amounts of wasted time and money, as well as missed deadlines."}
{"text": "Treats and chews are given to dogs for a variety of reasons. Chews are used to occupy the dog and help clean its teeth. Treats may have the same cleaning function, as well as provide some nutritional value.\nTreats and chews can be made from a variety of materials. For example, the products can be made of processed, dried animal skin, such as rawhide (untanned cattle skin) or pigskin. The skin is processed by removing hair and cleaning it of fat and muscle. The skin is then manipulated into the desired shape, such as by folding or rolling, and then dried.\nIn the prior art, treats and chews can be made of materials other than hide. For example, the products can be made of plastic, rubber, etc. \nEven when dogs chew on treats and chews, their breath may have an unpleasant odor and their teeth may not be cleaned. For example, on a rawhide chew, a dog has a hard time piercing the tough dried skin with its teeth. Thus, the teeth are not cleaned and the dog continues to have bad breath."}
{"text": "1. Technical Field\nThe present invention relates to computer and digital data compression, and more specifically to preventing rounding errors that can accumulate in MPEG-2 type decompression.\n2. Description of the Prior Art\nDigitized images require a large amount of storage space to store and a large amount of bandwidth to transmit. A single, relatively modest-sized image, having 480 by 640 pixels and a full-color resolution of twenty-four bits per pixel (three 8-bit bytes per pixel), occupies nearly a megabyte of data. At a resolution of 1024 by 768 pixels, a 24-bit color screen requires 2.3 megabytes of memory to represent. A 24-bit color picture of an 8.5 inch by 11 inch page, at 300 dots per inch, requires as much as twenty-five megabytes to represent.\nVideo images are even more data-intensive, since it is generally accepted that for high-quality consumer applications, images must occur at a rate of at least thirty frames per second. Current proposals for high-definition television (HDTV) call for as many as 1920 by 1035 or more pixels per frame, which translates to a data transmission rate of about 1.5 billion bits per second. This bandwidth requirement can be reduced somewhat if one uses 2:1 interleaving and 4:1 decimation for the xe2x80x9cUxe2x80x9d and xe2x80x9cVxe2x80x9d chrominance components, but 0.373 billion bits per second are still required.\nTraditional lossless techniques for compressing digital image and video information, such as Huffman encoding, run length encoding and the Lempel-Ziv-Welch algorithm, are far from adequate to meet this demand. For this reason, compression techniques which can involve some loss of information have been devised, including discrete cosine transform techniques, adaptive discrete cosine transform techniques, and wavelet transform techniques. Wavelet techniques are discussed in DeVore, Jawerth and Lucier, Image Compression Through Wavelet Transform Coding, IEEE Transactions on Information Theory, Vol. 38, No. 2, pp. 719-746 (1992); and in Antonini, Barlaud, Mathieu and Daubechies, Image Coding Using Wavelet Transform, IEEE Transactions on Image Processing, Vol. 1, No. 2, pp. 205-220 (1992).\nThe Joint Photographic Experts Group (JPEG) has promulgated a standard for still image compression, known as the JPEG standard, which involves a discrete cosine transform-based algorithm. The JPEG standard is described in a number of publications, including the following incorporated by reference herein: Wallace, The JPEG Still Picture Compression Standard, IEEE Transactions on Consumer Electronics, Vol. 38, No. 1, pp. xviii-xxxiv (1992); Purcell, The C-Cube CL550 JPEG Image Compression Processor, C-Cube Microsystems, Inc. (1992); and C-Cube Microsystems, JPEG Algorithm Overview (1992).\nAn encoder using the JPEG algorithm has four steps: linear transformation, quantization, run-length encoding (RLE), and Huffman coding. The decoder reverses these steps to reconstitute the image. For the linear transformation step, the image is divided up into 8*8 pixel blocks and a Discrete Cosine Transform is applied in both spatial dimensions for each block. The purpose of dividing the image into blocks is to overcome a deficiency of the discrete cosine transform algorithm, which is that the discrete cosine transform is seriously non-local. The imager is divided into blocks in order to overcome this non-locality by confining it to small regions, and doing separate transforms for each block. However, this compromise has a disadvantage of producing a tiled appearance (blockiness) upon high compression.\nThe quantization step is essential to reduce the amount of information to be transmitted, though it does cause loss of image information. Each transform component is quantized using a value selected from its position in each 8*8 block. This step has the convenient side effect of reducing the abundant small values to zero or other small numbers, which can require much less information to specify.\nThe run-length encoding step codes runs of same values, such as zeros, in items identifying the number of times to repeat a value, and the value to repeat. A single item like xe2x80x9ceight zerosxe2x80x9d requires less space torepresent than a string of eight zeros, for example. This step is justified by the abundance of zeros that usually result from the quantization step.\nHuffman coding translates each symbol from the run-length encoding step into a variable-length bit string that is chosen depending on how frequently the symbol occurs. That is, frequent symbols are coded with shorter codes than infrequent symbols. The coding can be done either from a preset table or one composed specifically for the image to minimize the total number of bits needed.\nSimilarly to JPEG, the Motion Pictures Experts Group (MPEG) has promulgated two standards for coding image sequences. The standards are known as MPEG-1 and MPEG-2. The MPEG algorithms exploit the common fact of relatively small variations from frame to frame. In the MPEG standards, a full image is compressed and transmitted only once for every twelve frames. The JPEG standard is typically used to compress these xe2x80x9creferencexe2x80x9d or xe2x80x9cintraxe2x80x9d frames. For the intermediate frames, a predicted frame is calculated and only the difference between the actual frame and the, predicted frame is compressed and transmitted. Any of several algorithms can be used to calculate a predicted frame, and the algorithm is chosen on a block-by-block basis depending on which predictor algorithm works best for the particular block. Motion detection can be used in some of the predictor algorithms. MPEG 1 is described in detail in International Standards Organization (ISO) CD 11172.\nAccordingly, for compression of video sequences the MPEG technique is one which treats the compression of reference frames substantially independently from the compression of intermediate frames between reference frames. The present invention relates primarily to the compression of still images and reference frames for video information, although aspects of the invention can be used to accomplish video compression even without treating reference frames and intermediate frames independently.\nThe above techniques for compressing digitized images represent only a few of the techniques that have been devised. However, none of the known techniques yet achieve compression ratios sufficient to support the huge still and video data storage and transmission requirements expected in the near future. The techniques also raise additional problems, apart from pure compression ratio issues. In particular, for real time, high-quality video image decompression, the decompression algorithm must be simple enough to be able to produce thirty frames of decompressed images per second. The speed requirement for compression is often not as extreme as for decompression, since for many purposes, images can be compressed in advance. Even then, however, compression time must be reasonable to achieve commercial objectives. In addition, many applications require real time compression as well as decompression, such as real time transmission of live events. Known image compression and decompression techniques which achieve high compression ratios, often do so only at the expense of requiring extensive computations either on compression or decompression, or both.\nThe MPEG-2 video compression standard is defined in ISO/IEC 13818-2 xe2x80x9cInformation technologyxe2x80x94Generic coding of moving pictures and associated audio information: Videoxe2x80x9d. MPEG-2 uses motion compensation on fixed sized rectangular blocks of pixel elements (xe2x80x9cmacroblocksxe2x80x9d) to use temporal locality for improved compression efficiency. The location of these xe2x80x9cmacroblocksxe2x80x9d in the reference pictures is given on half pixel boundaries, and so requires an interpolation of pixel elements. Such interpolation is specified in the MPEG-2 standard, as follows:\nThe xe2x80x9c//xe2x80x9d operator at the ends of cases B, C, and D is defined in the MPEG-2 specification as: xe2x80x9cInteger division with rounding to the nearest integer. Half-integer values are rounded away from zero unless otherwise specified [ . . . ]xe2x80x9d. Therefore, when a two or a four are the right hand operand and the left hand operand is greater or equal zero, the operator xe2x80x9c//xe2x80x9d can be replaced according to: x//2=(x+1) greater than  greater than 1; x//4=(x+2) greater than  greater than 2, where xe2x80x9c greater than  greater than xe2x80x9d denotes a binary right shift of the left hand operator by the right hand operator designated number of bits.\nSince xe2x80x9cavg(p00, p01):=(p00, p01)//2=(p00, p01+1) greater than  greater than 1xe2x80x9d is a simple operation to implement on a microprocessor, many commercial microprocessors include this operation as a built-in instruction. For example, xe2x80x9cpavgusbxe2x80x9d is included in Intel processors with the SSE extension, and AMD processors with the 3DNow! extensions. Such instruction is very useful in computing cases B and C, above.\nThe built-in xe2x80x9cpavgusbxe2x80x9d instruction is very tempting for use in implementing case-D, e.g., by executing it twice, as in xe2x80x9cavg(avg(p00, p01), avg(p10, p11))xe2x80x9d. This is also an interesting alternative for a pure hardware implementation, because the circuit to implement the first step of averaging can also be used a second time to do the second step. But to do so, will generate objectionable visual artifacts.\nThe usual sub-routine to implement case-D requires the execution of twenty-seven instructions, as is listed in the left column of Table I. But if xe2x80x9cpavgusbxe2x80x9d is used, the sub-routine can be reduced to five instructions, as is listed in the right column of Table I. So the temptation to use the quicker solution for case-D is very great.\nOne problem with the shortcut solution is that it causes, artifacts because rounding is handled different, than the MPEG standard dictates. So to have artifact-free implementations, the short-cut solution is prohibited.\nThe rounding error for a specific combination of four pixels can be calculated as:\nerr(p00, p01, p10, p11)=((((p00+p01+1) greater than  greater than 1)+((p10+p11+1) greater than  greater than 1)+1) greater than  greater than 1)xe2x88x92((p00, p01+p10+p11+2) greater than  greater than 2)\nThe overall error is only affected by the two least significant bits of each coefficient, so the total average error can be calculated as:   avgerr  =                                          xe2x80x83                    ⁢                                    xe2x80x83                                      i              =              0                        3                          ⁢                              xe2x80x83                    ⁢                                                    xe2x80x83                            ⁢                                                xe2x80x83                                                  j                  =                  0                                3                                      ⁢                          xe2x80x83                        ⁢                                          xe2x80x83                            ⁢                                                xe2x80x83                                                  k                  =                  0                                3                                      ⁢                          xe2x80x83                        ⁢                                          xe2x80x83                            ⁢                                                xe2x80x83                                                  l                  =                  0                                3                                      ⁢                          err              ⁢                              (                                  i                  ,                  j                  ,                  k                  ,                  l                                )                                                                4        ?        4        ?        4        ?        4              =                  3        8            .      \nThe rounding error is less than one least-significant-bit for first-generation compensation, and is not a problem. But the MPEG-2 standard allows for multi-generation compensation. The consequence of this is a predicted picture can be used as a reference picture for successive pictures. When a predicted picture is used as a reference picture, rounding errors can add up to more than one least significant bit, and the accumulated rounding errors can be large enough over several compensation generations to create objectionable visual artifacts.\nSuch visual artifacts or optical effects are called xe2x80x9cpumpingxe2x80x9d, and are perceived by the eye as image brightening over a series of predicted pictures that jump back to normal brightness at the next intra-picture. In particular, sequences that have an unusually high number of predicted pictures between intra codes pictures can experience serious pumping. For example, field structure encoded sequences or material that lack xe2x80x9cbi-directionalxe2x80x9d predicted pictures will suffer from pumping.\nThe invention comprises a method wherein two-step motion prediction errors for MPEG-2 interpolation case-D are corrected. An improved MPEG-2 decoder includes a logic gate, multiplexer, and adder. When both the horizontal (h0) and vertical (h1) motion vector components require a half pixel interpolation (case-D), the multiplexer forwards the constant minus three to the adder, otherwise a constant zero is used. Such adder modifies the D C coefficient input to the inverse discrete cosine transformer to include a correction term for the predicted pixels calculated by a two-step predictor. A correction value of xe2x88x920.375 is evenly distributed over all sixty-four resulting spatial coefficients during the inverse discrete cosine transform. This results statistically in a slightly brighter set of correction terms. Such offsets a slightly darker prediction that are formed by the two-step predictor. The output frames are statistically correct images."}
{"text": "1. Field\nAt least one example embodiment relates to a backlight unit, a display device including the backlight unit, and/or a method of manufacturing the backlight unit.\n2. Description of the Related Art\nIn general, a three-dimensional (3D) image is configured based on stereo view principles using two eyes of a person. A 3D image display may be a stereoscopic display and an autostereoscopic display. The autostereoscopic display acquires a 3D image by separating an image into a left image and a right image without using glasses, and uses, for example, a parallax barrier method and/or a lenticular method.\nHere, in the case of the parallax barrier method and/or the lenticular method, once a design is completed, the design may not be readily modified, optical efficiency may be low, and a crosstalk between a left image and a right image may be relatively high. In addition, some autostereoscopic displays employing the parallax barrier method and/or the lenticular method do not readily achieve a two-dimensional (2D) display. Alternatively, even if an autostereoscopic display employing the parallax barrier method and/or the lenticular method is capable of achieving the 2D display, the quality of a 2D image may be degraded."}
{"text": "This invention relates to a magnetic recording and optical reproducing apparatus in which a magneto-optical effect is utilized so as to reproduce information magnetically recorded on a recording medium by applying a recording magnetic field from a magnetic head.\nIt has been confirmed that an information bit length as short as 0.1 .mu.m in terms of a linear bit density can be recorded on a recording medium according to a recording method in which a recording magnetic field generated from a magnetic head is used to magnetically record information on the recording medium, as described in, for example, IEEE TRANSACTIONS ON MAGNETICS, VOL. MAG-23, No. 5, SEPTEMBER 1987, page 2070.\"\nOn the other hand, a method utilizing a magneto-optical effect for the high-sensitivity reproduction of information magnetically recorded on a recording medium is described in, for example, JP-A-55-153142."}
{"text": "1. Field of the Invention\nThe present invention relates to a composition for forming an antireflection film (an antireflection film-forming composition) suitable for microfabrication using a lithography process, in which a resist responding to various radiations such as a chemically amplified resist is used, a laminate with an antireflection film formed from the antireflection film-forming composition on a resist film, and a method for forming a resist pattern using the antireflection film-forming composition.\n2. Description of Background Art\nIn the field of microfabrication represented by fabrication of integrated circuit devices, the processing size has become more and more minute in order to achieve higher integration in a lithographic process. In recent years, development of a lithographic process enabling stable microfabrication with a line width of 0.5 μm or less has been actively proceeding.\nHowever, it is difficult to form such a fine pattern with high accuracy using a conventional method which utilizes visible rays (wavelength: 800-400 nm) or near ultraviolet rays (wavelength: 400-300 nm). To deal with this problem, a lithographic process using radiation with a shorter wavelength (wavelength: 300 nm or less) which can achieve a wider depth of focus and is effective for ensuring design rules with minimum dimensions has been proposed.\nAs a lithographic process using radiation with such a short wavelength, processes using deep ultraviolet rays such as a KrF excimer laser (wavelength: 248 nm) or an ArF excimer laser (wavelength: 193 nm), X-rays such as synchrotron radiation, and charged particle rays such as electron beams have been proposed. A chemically-amplified resist has attracted attention as a resist exhibiting high resolution for such short wavelength radiation. At present, development and improvement of the chemically amplified resist are important technical issues of the lithographic process.\nThe chemically-amplified resist contains a compound (a photoacid generator) which generates an acid upon irradiation (hereinafter called “exposure”). Chemical changes in the resist film (changes in polarity, breakage of a chemical bond, cross-linking reaction, etc.) caused by the catalytic action of an acid change solubility of the exposed area in a developer. A pattern is formed utilizing this phenomenon.\nA number of compositions for chemically-amplified resists, for example, a combination of a resin containing an alkali-soluble resin, in which the groups exhibiting affinity with an alkali are protected by a t-butyl ester group or t-butoxy carbonyl group, with an acid generator (for example, Japanese Patent Publication No. 59-45439), a combination of a resin containing an alkali-soluble resin, in which the groups exhibiting affinity with an alkali are protected by a silyl group, with an acid generator (for example, Japanese Patent Publication No. 60-52845), a combination of an acetal group-containing resin with an acid generator (for example, Japanese Patent Publication No. 2-25850), a combination of an alkali-soluble resin, a dissolution controller, and an acid generator, a combination of a novolac resin, a crosslinking agent, and an acid generator; and the like have been proposed.\nThese chemically-amplified resists, however, have a problem with storage stability, because the compositions are easily affected by moisture, oxygen, basic substances, and the like which are present in the atmosphere of the lithographic process. SPIE, Vol. 1466, “Advance in Resist Technology and Processing”, p. 2 (1991), for example, reports that a very small amount of dimethylaniline in the atmosphere deactivates acids existing near the surface of a resist film among the acids produced in the resist film by exposure, and forms a scarcely soluble layer on the surface of the film, which layer remains in the form of an eave-like projection on the surface of a resist pattern after development. Not only such a scarcely soluble layer decreases sensitivity and resolution of the resist, but also the eave-like projections formed in the resist pattern adversely affect the etching accuracy. The size of the eave-like projections tends to increase along with an increase of the time for which the resist is allowed to stand between each of the steps of exposure, post exposure baking, and development. This phenomenon is known as post exposure time delay (hereinafter referred to as “PED”), which unduly decreases the allowance of time in the lithographic process.\nAs a method for solving the problem of PED, a method of laminating an overcoat on the chemically-amplified resist film to block the film surface from the atmosphere has been proposed. For example, WO 92/05474 describes a method of laminating an overcoat of polyacrylic acid, polyvinyl butyral, polyvinyl alcohol, polystyrene sulfonic acid, or the like on the chemically-amplified resist film to inhibit invasion of basic materials into the resist film, thereby preventing a decrease of sensitivity and resolution of the resist. However, among such overcoats, polyacrylic acid, polyvinyl butyral, and polyvinyl alcohol cannot necessarily effectively prevent formation of the aforementioned scarcely-soluble layer, although they have a barrier effect. Polystyrene sulfonate, which has too strong acidity, has a drawback of initiating a chemical reaction due to the catalytic action of the acid in the chemically-amplified resist irrespective of exposure to light. Another problem of the overcoat, which is commonly applied to a resist film as an aqueous solution, is uneven coating due to insufficient wetting properties of the aqueous solution with the resist film.\nOn the other hand, since radiation commonly used in the lithographic process is light with a single wavelength, incidence radiation interferes with light reflected by the boundary surfaces on the top and bottom of the resist film. As a result, a phenomenon called “a standing-wave effect” or “a multiplex interference effect”, which is a phenomenon in which an effectual dose of radiation to which a resist film is exposed fluctuates due to mutual interference of radiations in the film according to fluctuation of the resist film thickness, irrespective of a constant actual dose of radiation, occurs and adversely affects formation of resist patterns. If the coating thickness varies due to slight difference of a resist composition, viscosity, resist film conditions, and the like, or the coating thickness fluctuates due to steps in the substrates (recessed parts are thicker than projected parts), such a film thickness difference changes the effectual dose of radiation to which a resist film is exposed, resulting in fluctuation of pattern dimensions and a decrease in accuracy of the resist pattern dimensions.\nTo overcome this problem caused by the standing-wave effect, a method of forming an antireflection film on the resist film to inhibit reflection on the coating surface and reduce multiplex interference in the film has been proposed. For example, a method of reducing the standing-wave effect by laminating layers of polysiloxane, poly(ethyl vinyl ether), polyvinyl alcohol, or the like as an antireflection film on a resist film is described in J. Electrochem. Soc., Vol. 137, No. 12, p. 3,900 (1990). In this instance, the reflex inhibition effect on the surface of the resist film mainly depends on the refractive index and film thickness of the antireflection film. The refractive index of an ideal antireflection film is the square root of n (n is the refractive index of the resist), and the thickness of an ideal antireflection film is an anisoploid of λ/4 m (wherein λ is the wavelength of radiation and m is the refractive index of the antireflection film).\nHowever, the antireflection film made from polysiloxane, poly(ethyl vinyl ether), or polyvinyl alcohol has a basic problem of being incapable of sufficiently inhibiting the standing-wave effect due to the small difference of the refractive index with that of a resist film. In addition, the antireflection film made from polysiloxane which is insoluble in water or a developer must have a separate step of removing the antireflection film before development using an antireflection film removing agent. The solubility in water or a developer of poly(ethyl vinyl ether) or polyvinyl alcohol is not necessarily sufficient. The antireflection film made from these polymers may leave residues on the resist film which impair resist performance such as resolution, developability, pattern profile, and the like. An additional problem of the overcoat which is applied to a resist film as an aqueous solution such as poly(ethyl vinyl ether) or polyvinyl alcohol is uneven coating due to insufficient wetting properties of the aqueous solution with the resist film.\nIn order to solve the problems in conventional antireflection films mentioned above, the inventor of the present invention has proposed (a) a copolymer of an acrylamide compound such as 2-acrylamide-2-methylpropane sulfonic acid in which the amide group has a sulfonic acid group bonded to the nitrogen atom via an organic group, and a fluoroalkyl acrylate compound such as 2,2,3,3,3-pentafluoropropyl acrylate and (b) a basic material blocking antireflection film containing a fluoroalkylsulfonic acid having an fluoroalkyl group with 5-15 carbon atoms and/or a fluoroalkylcarboxylic acid having an fluoroalkyl group with 5-15 carbon atoms, and a resist pattern forming method using the antireflection film (see Japanese Patent Application Laid-open No. 7-234514). This antireflection film can reduce the effect of a basic substance in the atmosphere and the standing-wave effect. However, the composition used for forming the antireflection film has a problem in coatability onto a resist film, which results in production of micro-bubbles in the order of micrometers in the formed coating.\nAn object of the present invention is to provide an antireflection film-forming composition having excellent coatability, capable of significantly inhibiting production of fine micro-bubbles and capable of forming an antireflection film with a sufficiently decreased standing-wave effect, and having excellent solubility in water and alkaline developers, a laminate with an antireflection film formed from the antireflection film-forming composition on a resist film, and a method for forming a resist pattern using the composition for forming a basic material blocking antireflection film."}
{"text": "A. Field of the Invention\nThe present invention relates to a new and improved game apparatus for setting up elongated objects, such as dominos, in predetermined configurations.\nB. Description of the Prior Art\nBoth children and adults enjoy setting up objects such as dominos in predetermined configurations and then through a chain reaction, causing the dominos to fall over sequentially and while falling over, to initiate an entertaining series of events. Not only are dominos set up in this manner for pleasure at home, but there are also contests and groups of people that get together to set up large chains of dominos over a large area and in many different configurations. Apparatus for use in conjunction with such play is described in U.S. Pat. Nos. 3,283,439, 3,283,440, 3,315,404, 3,621,601, 4,047,322 and 4,248,433.\nOne problem, however, is that often times an individual does not have the ability or the patience to place the dominos in very long, unusual and entertaining configurations. Often the tedium involved in setting up the dominos, an operation which is extremely delicate and often requires repeated efforts to set up an array without accidentally initiating a chain reaction, may exceed the potential enjoyment. Moreover, it would be desirable to have mechanical means for starting and signaling the finish of the chain reaction of falling dominos. It is also desirable to have several different devices in the path of the falling dominos to create different actions and then to continue the chain reaction of the falling dominos onto another chain of dominos. Such items are lacking in the prior art."}
{"text": "Electronically controlled machine tools and computer aided manufacturing software together may be utilized to create and/or form a manufactured component. This process often may include forming one or more features, such as holes, threaded holes, and/or slots, in a workpiece to transform the workpiece into at least a portion of the manufactured component.\nHistorically, programming an electronically controlled machine tool to form the feature has been a labor-intensive process in which the computer aided manufacturing software is utilized to describe each feature that is to be formed in the workpiece in detail. This process often becomes highly repetitive, especially when a plurality of similar features is to be formed in the workpiece.\nIn addition, newly designed cutting tool formations are continually being developed for more efficient processes for the formation of certain features. Current computer aided manufacturing software algorithms are not always capable of producing the most efficient or applicable tool path motion for a given cutting tool. The current computer aided manufacturing software algorithms also may cause significant and/or uneven wear of a cutting tool that is utilized to form the feature and/or may require that one or more manual machining operations be performed to ensure the integrity of the feature. As an example, machining a threaded hole with an electronically controlled machine tool may include drilling a hole into the workpiece and subsequently threading the hole to define the threaded hole. The machining process may utilize a combination cutting tool consisting of a drill, a thread mill, and a chamfer in one application. This combination tool is required to extend from an entrance of the hole into the hole to define the minor diameter of the threads within the hole and to chamfer the top of the hole at a specific surface feed rate. The combination tool then forms the threads of the threaded hole at a different surface feed rate. In such a machining operation, an endmost drilling portion of the tool removes a majority of the material, with a remainder of the combination tool, which includes a threading portion and a chamfering portion, removing less of the material. Thus, endmost thread milling cutters of the combination tool generally wear at an accelerated rate relative to a remainder of the combination tool. This may require that the combination tool be sharpened and/or replaced frequently and/or may require that the threaded hole be manually chased with a tap subsequent to formation of the threaded hole by the electronically controlled machine tool. Current tool path algorithms within computer aided manufacturing software do not have the capability to utilize such a combination tool in an optimal fashion and limit user control. Thus, there exists a need for associative templates, which may be defined by user inputs, for machining operations and/or for systems and/or methods that include and/or utilize the associative templates and thus may not limit the user to existing computer aided software algorithms and/or may decrease and/or eliminate repetitive tool motions when forming the feature."}
{"text": "The invention relates to a method and apparatus for applying liquids to plants. While the apparatus according to the present invention may be utilized to apply a wide variety of chemicals to plants, the apparatus and method are particularly suited for the application of liquid contact herbicides, such as ROUNDUP manufactured by Monsanto Company, to weeds.\nWeed control in growing crops is a continuing problem in agriculture, and a great deal of research has been done in that area and there have been many prior art proposals and procedures for the chemical treatment of weeds. Often times one particular structure is utilized for spraying herbicides in the crop rows between plants, and then another structure is utilized for contacting weeds taller than the crop plants.\nConventional overhead herbicide applicators include recirculating sprayers and rope wick applicators. Both structures can be successful if properly utilized, however recirculating sprayers are quite complex with a tendency to clog, and conventional rope wick applicators can apply less than the desired amount of herbicide to plants unless operating parameters are entirely appropriate. Further, neither system is readily adaptable to between row usage.\nAccording to the method and apparatus of the present invention, a rotating roller is utilized as a primary herbicide-applying component. While the utilization of rotating rollers per se has previously been proposed, such proposed structures conventionally have dripped herbicide onto rollers directly in contact with the ground and have been adapted to be rolled over the ground. Such structures therefore have limited utility for the application of herbicide to weeds adjacent crop rows. According to the present invention, a rotating liquid applicator roller is readily utilizable for applying herbicide simply and effectively to weeds between crop rows and to weeds growing above crop rows with equal efficiency. The apparatus according to the invention is readily adapted for mounting to a tractor tool bar, or to be pushed by hand.\nAccording to the method of the present invention, a liquid applicator roller is mounted for rotation about a generally horizontal axis and in operative communication with a reservoir of liquid herbicide. The method includes the steps of: Mounting the applicator roller above the ground and in a position for the peripheral surface thereof to be brought into contact with plants. Effecting movement of the applicator roller in a direction generally perpendicular to the axis of rotation of the roller. And, effecting powered rotation of the applicator roller so that it rotates into contact with plants and continuously transfers liquid from the surface thereof onto the plants, the surface being continuously replenished with liquid from the reservoir. Preferably rotation of the roller is effected in a direction of rotation so that it rotates into plants with which it comes into contact, providing a thorough wiping effect of herbicide onto the plants. The method may be practiced by moving the roller just above ground level between rows of crops, while positively preventing the roller from contacting crops in the rows on either side thereof; or by moving the roller above the tops of the crop plants to contact any plants extending upwardly above the tops of the crop plants.\nThe apparatus according to the invention includes a liquid applicator roller and a liquid reservoir and means for transferring liquid from the reservoir to the applicator roller surface. The applicator roller preferably comprises a tube of hard plastic having the ends thereof capped with spindles extending outwardly therefrom to define the axis of rotation of the roller. A housing mounts the roller for rotation about a given axis while covering the roller at the top and sides but being open at the bottom to allow the peripheral surface of the roller to rotate out of the housing. The housing also mounts the liquid transferring means. Means are provided for mounting the housing and the roller applicator above the ground, with the axis of rotation of the applicator roller generally horizontal. Preferably a drive wheel is provided, and means are provided for mounting the drive wheel to the housing so that it is rotatable about an axis parallel to the axis of rotation of the applicator roller, with surface means formed on the applicator roller and drive wheel effecting rotation of the applicator roller in response to the rotation of the drive wheel. Preferably the direction of rotation of the applicator roller is opposite to the direction of rotation of the drive wheel.\nWhen the drive wheel is utilized, it is mounted to one side of the housing while the other side of the housing is either mounted to the tractor tool bar or to a pair of ground engaging wheels. Horizontally spaced and elongated plates mounted to the housing adapt the housing for mounting either to a tool bar or to ground engaging wheels. Where a drive wheel is not utilized, the applicator roller may be powered by a motor mounted on a tractor.\nA liquid pool is defined within the housing substantially the entire width thereof, corresponding to the applicator roller length. The level of liquid in the pool is substantially even with the axis of rotation of the applicator roller. Transfer of liquid from the pool to the roller may be effected utilizing a sponge, or a pair of transfer rollers. Where the apparatus is utilized between crop rows, preferably a wiper is provided mounted to the housing to wipe any liquid off the roller periphery after plant contact and before a given portion of the roller rotates back into the housing through the housing open bottom. A more remote liquid reservoir may be provided, with a conduit leading from the liquid reservoir to the pool and with a level-responsive valve located in association with the conduit.\nIt is the primary object of the present invention to provide a method and apparatus for simply, efficiently, and effectively applying a liquid to plants, and particularly for applying a liquid herbicide. This and other objects of the invention will become clear from an inspection of the detailed description of the invention, and from the appended claims."}
{"text": "Businesses that offer services to large numbers of people often have access to large amounts of sensitive information. For example, many businesses perform thousands or millions of credit card transactions each year. Credit card information is just one example of sensitive information that many companies handle as part of providing a service. Often times, large amounts of data must be stored and accessed quickly as part of providing services involving sensitive information.\nCurrently existing methods for storing and accessing large amounts of sensitive information have disadvantages. Typically, when the information includes millions of records, the information is stored by a unique number, for example by consecutively numbered records This allows for easy storage, but accessing the information at a later date can be time consuming when there are large numbers of records, for example in the millions.\nIt is important for businesses to be able to quickly store and access any record of a large amount of records including sensitive information such as credit card information."}
{"text": "This invention relates to electric networks, and more particularly it relates to a phase sequence network or filter, responsive to selected symmetrical components of current in a polyphase alternating current electric power circuit.\nThe theory of symmetrical components is well understood by those skilled in the art of electric power transmission and distribution. In accordance with this theory, the phase currents or voltages in any unbalanced three-phase alternating-current electric circuit can be resolved into three sets of symmetrical, balanced current or voltages vectors known respectively as positive sequence, negative sequence and zero sequences components. The negative and zero sequence components of circuit current and voltage are not present under balanced circuit conditions. A more detailed discussion of the theory of symmetrical components can be found in the textbook, Symmetrical Components, by Wagner and Evans, published by McGraw-Hill in 1933.\nIt has been known that certain circuits, called sequence networks or filters, can be connected to an electric power system to provide an output signal that is proportional to the magnitude of the negative sequence components of current. These sequence networks or filters are particularly useful in the protective relaying art as they will sense the presence of negative sequence components of current which flow when the power system becomes unbalanced due to abnormal of fault conditions. It has also been known in the protective relaying art that it is desirable to have a sequence network that has a \"mixed\" response, that is, the output is a function of both the positive sequence components of current and the negative sequence components of current. One of the disadvantages of the prior art sequence networks that have a mixed response is that when the circuit values are adjusted to change the circuit response to the positive sequence components, the circuit's response to the negative sequence components also changes. I have found it desirable to have a sequence network in which the positive sequence response can be selectively varied without affecting the amplitude or phase of the circuit's response to the negative sequence components.\nThe output of the sequence sensitive circuit is a sine wave which generally must be further processed to derive a d-c control signal. Since the processing of the sine wave signal usually involves some filtering which delays the detection of the fault condition, it is desirable to keep the filter as small as possible.\nIt is, therefore, one object of this invention to provide a sequence network, responsive to both the positive sequence components and negative sequence components of a three phase alternating current power system, in which the response to the positive sequence components can be selectively varied without affecting the response to the negative sequence components.\nIt is another object of this invention to provide a mixed sequence responsive circuit in which the positive sequence response is 180.degree. out of phase with the negative sequence response.\nAnother object of this invention is to reduce the filtering required to derive a control signal from the sine wave output of a sequence sensitive circuit."}
{"text": "The sport of towing a rider on a water planing device behind a boat is well known. The types of devices may generally be divided into the categories of water skis and water planes or aquaplanes. In the categories of water skis the advantages of control and versatility are available. However, water skiing requires a development of sufficient strength to hold onto a tow rope while maintaining the skis in a parallel relationship and standing erect while doing both. This makes waterskiing a difficult and often frustrating sport for the beginner, requiring the simultaneous development of arm strength, leg strength and balance.\nWater planes are generally a single flat, wide board on which the rider stands. In some versions, for example the water sports vehicle disclosed in U.S. Pat. No. 3,581,328 of Smith, the problem of direct force on the rider from the tow rope is alleviated by connecting the tow rope to the front of the vehicle and having a handle extending from the front for the rider to hold for stability without requiring a great deal of upper body strength. Typically, however the water planing devices fail to provide as much control as is available with water skis, requiring considerable skill to negotiate turns while remaining erect.\nSeveral devices have been disclosed which are hybrids between the two aforementioned categories including the aquaplane of U.S. Pat. No. 2,841,805 issued to Roudebush and the U-shaped water ski described in U.S. Pat. No. 3,585,664 of Thompson. These devices primarily take parallel skis and join them by a rigid or semirigid crossbar at their front, thereby maintaining the skis in a permanent parallel relationship. Similar to water skis, these hybrid devices require a relatively high minimum speed to prevent the rider from sinking, further requiring a substantial amount of resistance to be overcome, thus considerable leg strength and balance, in order to initiate planing. Further, because of the fixed location of the tow rope attachment at the front of the device, the only possible means to adjust for different riders is variation of the length of the rope extension from the front of the device to the handles onto which the rider holds. No adjustment is provided for either the rider's weight or skill level, or for allowing the device to plane at different speeds.\nIt would be desirable to provide a water ski board for towing behind a boat or the like which may be used by beginners or accomplished skiers, by providing the stability needed by beginners as well as adjustability for speed and skill level for increasing the challenge and controllability for advanced skiers. It is to such a device that the present invention is directed."}
{"text": "1. Field of the Invention\nThe present invention relates to a clutch assembly, and more particularly to a frictional clutch assembly that can transmit power in a unidirectional manner.\n2. Description of Related Art\nWith reference to FIG. 6, a conventional unidirectional transmitting assembly (10) in accordance with the prior art comprises a driving axle (11), a transmitting sleeve (13), multiple rollers (12) and an output axle (14). The transmitting sleeve (13) is received in the driving axle (11) and has multiple cutouts (15) defined in the outer surface of the sleeve (13). Each cutout (15) has a wide segment and a narrow segment and receives one of the rollers (12). The output axle (14) is mounted in and engages with the transmitting sleeve (13) to rotate with the transmitting sleeve (13).\nWhen the driving axle (11) rotates along a direction corresponding to the narrow segments of the cutouts (15), clockwise as shown by the black arrow, the rollers (12) will be pushed to the narrow segments of the cutouts (15) by means of the friction between the driving axle (11) and the rollers (12). Consequently, the rollers (12) will be securely engaged between the driving axle (11) and the transmitting sleeve (13), such that the transmitting sleeve (13) will rotate with the driving axle (11) and the output axle (14) will be driven to rotate.\nWhen the driving axle (11) rotates in an opposite direction that corresponds to the wide segments of the cutouts (15), counterclockwise as shown by the white arrow, the rollers (12) will move to the wide segments of the cutouts (15). Then, the rollers (12) will freely rotate in the cutouts (15), and the transmitting sleeve (13) and the output axle (14) will not rotate and will keep stationary. Accordingly, the input power will be transmitted in a unidirectional manner.\nHowever, the conventional unidirectional transmitting assembly (10) does not have a function of a clutch, so it is necessary to add an additional clutch device for a machine with a conventional unidirectional transmitting assembly and this renders the machine expensive.\nTo overcome the shortcoming, the present invention tends to provide a frictional clutch assembly to mitigate or obviate the aforementioned problems."}
{"text": "Aqueous is a clear, colorless fluid that fills the anterior and posterior chambers of the eye. The aqueous is formed by the ciliary body in the eye and is a carrier of nutrients for the lens. In addition, the aqueous provides a continuous stream into which surrounding tissues can discharge the waste products of metabolism.\nThe aqueous produced in the ciliary processes circulates from the posterior chamber to the anterior chamber of the eye through the pupil and is absorbed through the trabecular meshwork, a plurality of criss-crossing collagen cords covered by endothelium. Once through the trabecular meshwork, the aqueous passes through Schlemm's canal and into venous circulation. The rate of aqueous outflow through the trabecular meshwork in a normal eye is typically 2.1 .mu.L/min. Intraocular pressure in the eye is maintained by the formation and drainage of the aqueous. All the tissues within the corneoscleral coat covering the eyeball are subject to this pressure, which is higher than pressure exerted on tissues at other locations in the body.\nGlaucoma is a progressive disease of the eye characterized by a gradual increase of intraocular pressure. This increase in pressure is most commonly caused by stenosis or blockage of the aqueous outflow channel, resulting in excessive buildup of aqueous fluid in the eyeball. Other causes include increase in venous pressure outside the eye which is reflected back through the aqueous drainage channels and increased production of aqueous. In a \"normal\" eye, intraocular pressure ranges from 4 to 12 mm mercury. In an eye with glaucoma, this pressure can rise to as much as 50 mm mercury. This increase in intraocular pressure produces gradual and permanent loss of vision in the afflicted eye.\nExisting corrective methods for the treatment of glaucoma include drugs, surgery, and implants. Miotic drugs lower intraocular pressure by facilitating aqueous outflow. Beta blockers, epinephrine products, and carbonic anhydrase inhibitors which inhibit production of the aqueous, are also commonly used in pharmacological glaucoma treatment. Steroids have been used in long-term glaucoma treatment as well. However, pharmacological treatment is prohibitively expensive to a large majority of glaucoma patients. In addition, many people afflicted with the disease live in remote or undeveloped areas where the drugs are not readily accessible. The drugs used in the treatment, in particular the steroids, often have undesirable side effects and many of the long-term effects resulting from prolonged use are not yet known.\nSurgical procedures have been developed in an effort to treat victims of glaucoma. An iridectomy, removal of a portion of the iris, is often used in angle-closure glaucoma wherein there is an occlusion of the trabecular meshwork by iris contact. Removal of a piece of the iris then gives the aqueous free passage from the posterior to the anterior chambers in the eye. A trabeculotomy, opening the inner wall of Schlemm's canal, is often performed in cases of open-angle glaucoma so as to increase the outflow of the aqueous, thereby decreasing intraocular pressure. While often successful, these surgical techniques possess inherent risks associated with invasive surgery on an already afflicted eye. Furthermore, the tissue of the eye can grow back to the preoperative condition, thereby necessitating the need for further treatment.\nOcular implants are often used in long-term glaucoma treatment without the disadvantages of drugs and invasive surgery. One such implant is disclosed in U.S. Pat. No. 4,457,757 entitled \"Device for Draining Aqueous Humor\" and commercially available as the Molteno.TM. Seton Implant. The implant comprises a drainage tube connected to one or more ridged plate reservoirs. The reservoir plates are designed to conform to the curvature of the eye. A reservoir plate is placed under Tenon's capsule and sutured to the sclera. The drainage tube is implanted into the anterior chamber through a scleral flap. A second plate can be implanted under the superior rectus muscle and sutured to the sclera. At this point, the body will form a tissue around these plates. Increased pressure causes the tissues above the plates to lift and form a bleb into which aqueous fluid from the anterior chamber drains via the drainage tube. Once inside the bleb, the aqueous seeps into intercellular spaces and is removed by surrounding capillaries or lymphatics. This type of implant is disadvantageous as the plates are formed of a rigid plastic which makes insertion beneath the eye tissue difficult and time-consuming. Furthermore, the rigid material poses a risk of irritation and/or damage to adjacent vasculature and tissue.\nU.K. Patent Application 2,160,778 entitled \"Aqueous humor drainage device\" discloses a similar type of implant device comprising a drainage tube and a drainage body. The tube is fixed to and opens directly onto a surface of the body. The device is sutured to the sclera of the eye and the tube positioned within the anterior chamber so as to provide outflow for the aqueous contained therein. The device further includes a pressure gradient limiting valve formed as a slit in the tube, however, this type of valve does not allow patent, i.e., open or two-way, flow through the drainage tube, thereby preventing retrograde aqueous flow into the anterior chamber.\nGlaucoma implants require formation of scar tissue and a drainage bleb around the implant to control the outflow of the aqueous. In some cases, due to extreme intraocular pressure, the drainage bleb that forms is excessive in size, and may cause impaired vision due to unexpected complications. If the drainage bleb is too large, the eye may protrude from its orbit, or eye socket, and cause distortions in the patient's vision. This is especially problematic in patients with small orbits. In addition, the drainage bleb may press upon the internal tissues of the eye which will result in impaired vision due to increased retinal pressure. In these cases, the potential side effects may negate the advantages of the surgery."}
{"text": "The present invention, in some embodiments thereof, relates to methods and compositions for treating angiogenesis-related diseases such as cancer, and methods of selection thereof.\nThe neuropilin-1 (np1) and the neuropilin-2 (np2) receptors were originally characterized as functional receptors for axon guidance factors belonging to the class-3 semaphorin (sema3) family. It was subsequently realized that the neuropilins are expressed by endothelial cells and by many types of cancer cells. It was also found that the neuropilins function in addition as receptors for several angiogenic factors belonging to the VEGF family and as receptors for the angiogenesis/metastasis inducing growth factor hepatocyte growth factor/scatter factor (HGF/SF), and that they function as potent enhancers of their pro-angiogenic activity.\nMost of the sema3s, with the exception of sema3E which binds to PlexD1, bind to one of the two neuropilin receptors or to both. Neuropilins form spontaneous complexes with several members of the plexin receptor family. In these complexes the sema3s bind to neuropilins while the plexins function as the signal transducing elements. The four type-A plexins (plexins-A1 to plexin-A4) as well as plexin-D1 form complexes with neuropilins and participate in neuropilin mediated signal transduction.\nSemaphorins sema3B and sema3F were also characterized as tumor suppressors whose loss contributes to the development of lung cancer [Tomizawa, Y., 2001, Proc. Natl. Acad. Sci. U.S.A 98:13954-13959; Xiang, R., 2002. Cancer Res. 62:2637-2643].\nThe identification of neuropilins in endothelial cells suggested that class-3 semaphorins may be able to regulate angiogenesis. Indeed, the class-3 semaphorin sema3F, a np2 agonist, functions as a repellant of endothelial cells, induces apoptosis of endothelial cells upon prolonged stimulation [Bielenberg, D. R., et al., 2004, J. Clin. Invest 114:1260-1271; Guttmann-Raviv, N., et al., 2007, J. Biol. Chem. 282:26294-26305] and inhibits angiogenesis and tumor progression in-vivo [Bielenberg, D. R., et al., 2004, J. Clin. Invest 114:1260-1271; Kessler, O., et al 2004. Cancer Res. 64:1008-1015. The np1 agonist sema3A was also shown to inhibit in-vitro and in-vivo angiogenesis [Miao, H. Q., 1999. J. Cell Biol. 146:233-242. Bates, D., 2003, Dev. Biol. 255:77-98; Acevedo, L. M., 2008. Blood. 111:2674-2680]. Sema3E was also characterized as a repulsive agent that inhibits the invasion of PlexD1 expressing blood vessels into somites during embryonic development [Gu C. et al., 2005, Science 307:265-268].\nIn contrast, existing data suggests that sema3C functions as a pro-tumorigenic and pro-angiogenic agent [Herman J G, et al, Int. J. Oncol. 2007; 30:1231-1238; Banu N, FASEB J. 2006; 20:2150-2152].\nThe fact that neuropilins and plexins such as PlexD1 are also expressed by many types of tumor cells indicates that semaphorins may also affect the behavior of tumor cells directly. Indeed, sema3s such as sema3F and sema3B have been observed to inhibit the adhesion, migration or the proliferation of tumor cells expressing appropriate semaphorin receptors [Tomizawa, Y., 2001, Proc. Natl. Acad. Sci. U.S. A 98:13954-13959; Xiang, R., 2002. Cancer Res. 62:2637-2643; Bielenberg, D. R., et al., 2004, J. Clin. Invest 114:1260-1271; Nasarre, 2006, Neoplasia. 7:180-189]. In contrast, however, the cleavage product of Sema3E, was shown to be an inducer of tumor invasiveness and tumor metastasis [Christensen, C, 2005, Cancer Res. 65, 6167-6177]."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory capable of executing a read test for stored contents comparatively easily, and more particularly to a nonvolatile semiconductor memory such as an electrically writable/erasable flash memory.\n2. Description of the Background Art\nIn general, the production cost of a memory implies the total of the cost required for a wafer process, an assembly and a test. The cost of the test depends on how many chips can be tested per unit time by means of one tester. In order to reduce the test cost to produce a more inexpensive memory, accordingly, it is required that the prolongation of a test time caused by an increase in a storage capacity of the memory should be minimized even if the storage capacity of the memory is increased twofold to fourfold with an increase in the capacity.\nIn order to shorten a test time of the memory, the following should be implemented:\n(1) a reduction in an operation time required for write, read or the like;\n(2) development of a test pattern having a higher defect detecting capability; and\n(3) development of a test mode capable of carrying out write/read at a higher speed.\nReferring to (1), an increase in a speed of production has been required. Therefore, the test time tends to be shortened comparatively easily without a special contrivance through an enhancement in transistor performance by microfabrication and a reduction in a load capacity.\nReferring to (2), there are various test patterns having high defect detecting capabilities. Typically, a checker board pattern can be taken as an example.\nFIG. 95 is a diagram illustrating an example of the checker board pattern. FIG. 95 shows a checker board pattern having one bit/cell (binary-value).\nAs shown in FIG. 95, a checker board pattern CHK2 is a test pattern having a repetitive cycle of 2 bits in which adjacent bits always have a relationship of xe2x80x9c0xe2x80x9d and xe2x80x9c1xe2x80x9d. The checker board pattern CHK2 can detect open (disconnection) and a short circuit of a word line, open and a short circuit of a bit line, and a defective short circuit of floating gates in a nonvolatile semiconductor memory represented by a flash memory.\nFIG. 96 is a diagram illustrating a checker board pattern CHK4 having 2 bits/cell (quaternary-value). As shown in FIG. 96, 2-bit patterns CHK4-A to CHK4-D are repeated every 4-bit repetitive cycle in order to correspond to the quaternary-value (2 bits/cell) to be multiple-valued storage through application of the checker board pattern CHK2 shown in FIG. 95.\nFIG. 97 is a diagram illustrating a checker board pattern CHK8 having 3 bits/cell (octal-value). As shown in FIG. 97, a 3-bit pattern is repeated every 8-bit repetitive cycle in order to correspond to the octal-value (3 bits/cell) to be multiple-valued storage.\nIn the checker board pattern CHK4 and the checker board pattern CHK8, it is possible to detect a defective mode of the checker board pattern CHK2, and furthermore, to detect that all multiple-valued data can be written and read or not in the same word line and the same bit line.\nReferring to (3), there have been various methods. In general, if a test mode is incorporated into a chip, a chip area tends to be increased due to an incorporated circuit. Consequently, a cost required for a wafer process is increased. Accordingly, in the case in which the test mode is to be incorporated into the chip, it is necessary to note that the total production cost should be minimized.\nNext, the trend of a product of a flash memory and that of development will be described.\nAs an alternative to an EPROM (electrically writable nonvolatile semiconductor memory), a flash memory has spread for code storage. In recent years, a flash memory for mass data storage has spread more increasingly than the flash memory for code storage. In the case in which the flash memory for data storage carries out writing and reading randomly, it operates at a lower speed than that of the flash memory for code storage. However, in the case in which the flash memory for data storage carries out writing and reading sequentially, it can operate at a higher speed than that of the flash memory for code storage. The flash memory for data storage has had a larger capacity to exceed a DRAM through microfabrication of a processing pattern and a multiple-valued technique for storing multibit data in one memory cell.\nNext, a test time for the flash memory will be described.\nThe flash memory for code storage has such a structure that a reading operation can be carried out at a much higher speed than a writing operation. Therefore, a time required for a read test can be almost ignored with respect to the whole test time. However, the flash memory for data storage carries out the writing operation at a higher speed and has a larger capacity than the flash memory for code storage. Therefore, the time required for the read test cannot be disregarded with respect to the whole test time.\nFor example, a 256 Mbit flash memory has a 16 Ksector structure in which a 2 Kbyte write/read unit (hereinafter referred to as a xe2x80x9csectorxe2x80x9d) is present for 16K. Approximately 50 xcexcs is required for a reading head every sector (hereinafter referred to as a xe2x80x9c1st accessxe2x80x9d) and 50 ns is required for subsequent data transfer every byte.\nAccordingly, approximately 2.5 s ((50 xcexcs+50 nsxc3x972 Kbyte)xc2x716 Ksector) is required for carrying out a read test in the whole area of the 256 Mbit flash memory.\nIn a probing check stage of a wafer state, furthermore, it is hard to carry out a read test with a 2nd access 50 ns of a product specification due to a resistance and a capacitance of a probe needle and a resistance and a capacitance of a probe card, and a test time is further increased.\nIn general, the flash memory has an automatic writing/erasing function. The automatic writing function implies a function of repeating a write pulse applying operation and an operation (hereinafter referred to as a xe2x80x9cverify operationxe2x80x9d) for deciding whether desirable data are written (or erased) to (or from) an object memory cell for writing in accordance with a logic circuit (hereinafter referred to as a xe2x80x9ccontrol circuitxe2x80x9d) provided in an EEPROM, ending the repetition of the write pulse applying operation and the verify operation when it is decided that all the object memory cells store the desirable data and outputting a signal for giving, to the outside of the EEPROM, a notice that the writing operation (or the erasing operation) has been completed.\nIn order to decide whether xe2x80x9cthe contents stored in all the object memory cells are the desirable dataxe2x80x9d, an all latch deciding circuit (hereinafter referred to as an xe2x80x9cALL deciding circuitxe2x80x9d) is provided. The ALL deciding circuit serves to decide that all sense latches in a sense latch group provided for storing the result of read of the memory cell are xe2x80x9c1xe2x80x9d or xe2x80x9c0xe2x80x9d.\nFIG. 98 is a block diagram schematically showing a conventional ALL deciding circuit and a periphery thereof. In this specification, it is assumed that the case in which xe2x80x9c1xe2x80x9d is written to the memory cell is set to xe2x80x9cwritexe2x80x9d and the case in which xe2x80x9c0xe2x80x9d is written to the memory cell is set to xe2x80x9cerasexe2x80x9d. Referring to FIG. 98, description will be given in which a left memory cell group 31 is referred to as an L mat 31 and a right memory cell group 32 is referred to as an R mat 32.\nAs shown in FIG. 98, a sense latch group 33 is provided between the L mat 31 and the R mat 32. The sense latch group 33 transmits and receives data in a sector unit to and from the L mat 31 or the R mat 32. The latch data of the sense latch group 33 are output to an ALL deciding circuit 34.\nThe ALL deciding circuit 34 receives control signals LorR, 0or1 and ENABLE from an external control CPU 35 and outputs a decision result ALL34 to the control CPU 35. xe2x80x9c0xe2x80x9d/xe2x80x9c1xe2x80x9d of the LorR designates reading from the L mat 31/R mat 32, xe2x80x9c0xe2x80x9d/xe2x80x9c1xe2x80x9d of the 0or1 designates write verify/erase verify, and xe2x80x9c0xe2x80x9d/xe2x80x9c1xe2x80x9d of the ENABLE designates inactivity/activity of the ALL deciding circuit 34.\nFIGS. 99 to 102 are diagrams illustrating an operation principle of a conventional ALL deciding circuit, in which the ALL deciding circuit has a function of deciding that all sense latches are xe2x80x9c0xe2x80x9d.\nAs shown in FIGS. 99 to 102, the sense latch group 33 is provided between the L mat 31 and the R mat 32, and nodes N11 and N12 of each latch L33 of the sense latch group 33 are connected to the L mat 31 and the R mat 32 in one memory cell unit, respectively. The node N11 and the node N12 in the latch L33 are constituted to have a logical inverting relationship. A data latch group 36 is provided on the opposite side of the sense latch group 33 with the L mat 31 interposed therebetween, and a data latch group 37 is provided on the opposite side of the sense latch group 33 with the R mat 32 interposed therebetween.\nFirst of all, the ENABLE is set to xe2x80x9c1xe2x80x9d for initialization and the ALL deciding circuit 34 is thus brought into an active state.\nAs shown in FIG. 99, when the write verify for the L mat 31 is to be carried out, LorR=xe2x80x9c0xe2x80x9d (L) and 0or1=xe2x80x9c1xe2x80x9d are set so that data read from the L mat 31 are latched onto each latch L33 of the sense latch group 33. If xe2x80x9c1xe2x80x9d is normally written to the L mat 31, the nodes N11 of all the latches L33 of the sense latch group 33 are set to xe2x80x9c1xe2x80x9d and the nodes N12 are set to xe2x80x9c0xe2x80x9d.\nAccordingly, it is possible to execute the L mat write verify by deciding the state of the nodes N12 of all the latches L33 in the sense latch group 33 (whether all of them are xe2x80x9c0xe2x80x9d) through the ALL deciding circuit (R side sense latch decision).\nAs shown in FIG. 100, similarly, when the erase verify for the L mat 31 is to be carried out, LorR=xe2x80x9c0xe2x80x9d and 0or1=xe2x80x9c0xe2x80x9d are set. If xe2x80x9c0xe2x80x9d is normally written to the L mat 31, the nodes N11 of all the latches L33 in the sense latch group 33 are set to xe2x80x9c0xe2x80x9d and the nodes N12 are set to xe2x80x9c1xe2x80x9d.\nAccordingly, it is possible to execute the L mat erase verify by deciding the state of the nodes N11 of all the latches L33 in the sense latch group 33 through the ALL deciding circuit (L side sense latch decision).\nAs shown in FIG. 101, when the write verify for the R mat 32 is to be carried out, LorR=xe2x80x9c1xe2x80x9d (R) and 0or1=xe2x80x9c1xe2x80x9d are set so that data read from the R mat 32 are latched onto each latch L33 of the sense latch group 33. If xe2x80x9c1xe2x80x9d is normally written to the R mat 32, the nodes N12 of all the latches L33 of the sense latch group 33 are set to xe2x80x9c1xe2x80x9d and the nodes N11 are set to xe2x80x9c0xe2x80x9d.\nAccordingly, it is possible to execute the R mat write verify by deciding the state of the nodes N11 of all the latches L33 in the sense latch group 33 through the ALL deciding circuit (L side sense latch decision).\nAs shown in FIG. 102, similarly, when the erase verify for the R mat 32 is to be carried out, LorR=xe2x80x9c1xe2x80x9d and 0or1=xe2x80x9c0xe2x80x9d are set. If xe2x80x9c0xe2x80x9d is normally written to the R mat 32, the nodes N12 of all the latches 33 in the sense latch group 33 are set to xe2x80x9c0xe2x80x9d and the nodes N11 are set to xe2x80x9c1xe2x80x9d.\nAccordingly, it is possible to execute the R mat erase verify by deciding the state of the nodes N12 of all the latches L33 in the sense latch group 33 through the ALL deciding circuit (R side sense latch decision).\nThe ALL deciding circuit carries out the R or L side sense latch decision based on a logical expression of {(LorR) X OR (0or1)}. It is sufficient that the R side sense latch decision is carried out with xe2x80x9c1xe2x80x9d and the L side sense latch decision is carried out with xe2x80x9c0xe2x80x9d.\nBased on the result of the decision of the ALL deciding circuit for deciding whether all the memory cells in the sector are xe2x80x9c0xe2x80x9d, thus, it is possible to execute the write verify and erase verify operations without outputting read data for each bit from data input/output pins.\nIn this case, a time of approximately 150 xcexcs (50 xcexcs+50 nsxc2x72 Kbyte) required for normal reading per sector can be shortened to 50 xcexcs+xcex1 (xcex1 less than 1 xcexcs). Therefore, the read test time can be shortened to approximately one-third.\nIn this case, however, it is premised that all the write data in the sectors are identical. For this reason, there has been a problem in that the write data cannot be used for test reading through the checker board pattern CHK2, CHK4, CHK8 or the like which has a high defect detecting capability.\nA first aspect of the present invention is directed to a semiconductor memory comprising a plurality of memory cells, each of which can store N-valued (Nxe2x89xa72) information, a data reading device for reading a predetermined number of read data from a predetermined number of memory cells out of said plurality of memory cells during a reading operation for a test, and a deciding device for classifying the predetermined number of read data into K (Kxe2x89xa72) groups and outputting a decision result based on whether all the read data in the respective K groups are identical during the reading operation for a test.\nA second aspect of the present invention is directed to the semiconductor memory according to the first aspect of the present invention, wherein the K includes N, the memory cells include memory cells arranged in a matrix defined by first and second directions, the predetermined number of memory cells include memory cells provided in the same position in the second direction and provided in series in the first direction, and the deciding device classifies the predetermined number of read data such that the predetermined number of memory cells are classified into the same groups at N intervals in the second direction.\nA third aspect of the present invention is directed to the semiconductor memory according to the second aspect of the present invention, wherein the N-value includes a 2m-value (mxe2x89xa71).\nA fourth aspect of the present invention is directed to the semiconductor memory according to any of the first to third aspects of the present invention, the deciding device includes a sense storing device for sensing and storing the predetermined number of read data, and a decision result output device for deciding whether all the read data in the respective K groups are identical based on stored contents of the sense storing device and for outputting a result of the decision.\nA fifth aspect of the present invention is directed to the semiconductor memory according to any of the first to fourth aspects of the present invention, wherein the N-value includes a multiple-value to be ternary or more, the reading operation for a test includes first to Lth (Lxe2x89xa72) partial reading operations for a test which have different reading conditions and the read data include first to Lth 1-bit read data, the data reading device reads the predetermined number of first to Lth 1-bit read data during execution of the first to Lth partial reading operations for a test, the result of decision includes first to Lth partial decision results, and the deciding device outputs an ith (i=1 to L) partial decision result based on whether all ith 1-bit read data in the respective K groups are identical during an ith partial reading operation for a test.\nA sixth aspect of the present invention is directed to a semiconductor memory comprising a plurality of memory cells, each of which can store N-valued (Nxe2x89xa72) information, a data reading device for reading a predetermined number of read data from a predetermined number of memory cells out of said plurality of memory cells during a reading operation for a test, an expectation storing device for storing a predetermined number of expectation data, and a deciding device for outputting a result of decision based on a result of comparison of the predetermined number of read data with the predetermined number of expectation data during the reading operation for a test.\nA seventh aspect of the present invention is directed to the semiconductor memory according to the sixth aspect of the present invention, the deciding device includes a sense storing device for sensing and storing the predetermined number of read data, and a decision result output device for outputting the result of decision based on a result of comparison of stored contents of the sense storing device with stored contents of the expectation storing device.\nAn eighth aspect of the present invention is directed to the semiconductor memory according to the seventh aspect of the present invention, wherein the N-value includes a multiple-value to be ternary or more, the reading operation for a test includes first to Lth (Lxe2x89xa72) partial reading operations for a test which have different reading conditions, the read data include first to Lth 1-bit read data and the expectation data include first to Lth 1-bit expectation data, the data reading device reads the predetermined number of first to Lth 1-bit read data every execution of the first to Lth partial reading operations for a test respectively, the result of decision includes first to Lth partial decision results, the deciding device outputs an ith (i=1 to L) partial decision result based on a result of comparison of the predetermined number of ith 1-bit read data with the predetermined number of ith 1-bit expectation data during the ith partial reading operation for a test, and the predetermined number of second to Lth 1-bit expectation data are obtained by changing the predetermined number of first to (Lxe2x88x921)th 1-bit expectation data based on the stored contents of the expectation storing device and the sense storing device, respectively.\nA ninth aspect of the present invention is directed to the semiconductor memory according to any of the sixth to eighth aspects of the present invention, wherein the expectation storing device includes a data storing device for temporarily storing data when transmitting and receiving data between the memory cells and an outside.\nAs described above, according to the first aspect of the present invention, a comparatively complicated test pattern having the same value set to the K groups is written to the memory cells and the reading operation for a test is then executed to obtain the result of decision. Consequently, the read test for the memory cells can be carried out at a high speed.\nMoreover, it is sufficient that the deciding device has the function of deciding whether all the read data in the respective K groups are identical. Therefore, the circuit area of the semiconductor memory is rarely increased due to the provision of the deciding device therein.\nAccording to the second aspect of the present invention, a predetermined number of memory cells are classified into the same group at the N intervals in the second direction. Therefore, it is possible to carry out the read test in which the checker board pattern having the N-value is set to be a test pattern.\nAccording to the third aspect of the present invention, it is possible to carry out a read test in which a checker board pattern having a repetitive cycle of m bits is set to be the test pattern.\nAccording to the fourth aspect of the present invention, a predetermined number of read data are sensed and stored in the sense storing device. Consequently, it is possible to obtain a result of decision with high precision.\nAccording to the fifth aspect of the present invention, the result of decision comprising the first to Lth partial decision results is obtained. Consequently, it is possible to carry out the read test in a multiple-valued storage state for the memory cell without hindrance.\nAccording to the sixth aspect of the present invention, a predetermined number of expectation data are stored from the outside into the expectation storing device. Consequently, it is possible to carry out a read test based on an optional test pattern. Moreover, in the case in which the predetermined number of expectation data are utilized in common between the predetermined number of data read plural times, it is preferable that the predetermined number of expectation data should be stored in the expectation storing device at a first time. Therefore, a time required for storing the predetermined number of expectation data in the expectation storing device can be omitted during the reading operation of the predetermined number of read data at and after a second time. Correspondingly, the read test can be carried out at a high speed.\nMoreover, it is preferable that the deciding device should have the function of outputting the result of decision based on the result of comparison of the predetermined number of read data with the predetermined number of expectation data. Therefore, the circuit area of the semiconductor memory is rarely increased due to the deciding device constituted therein.\nAccording to the seventh aspect of the present invention, the predetermined number of read data are sensed and stored in the sense storing device. Consequently, it is possible to obtain a result of decision with high precision.\nAccording to the eighth aspect of the present invention, the result of decision comprising the first to Lth partial decision results is obtained. Consequently, it is possible to carry out the read test in a multiple-valued storage state for the memory cell without hindrance. In this case, the predetermined number of second to Lth 1-bit expectation data are obtained by changing the predetermined number of first to (Lxe2x88x921)th 1-bit expectation data based on the contents stored in the expectation storing device and the sense storing device, respectively. Therefore, it is sufficient that only the predetermined number of first 1-bit expectation data should be stored in the expectation storing device.\nAccording to the ninth aspect of the present invention, the expectation storing device does not need to be added specially for expectation data storage.\nAn object of the present invention is to provide a semiconductor memory capable of executing a read test at a high speed based on a comparatively complicated test pattern without increasing a circuit area. These and other objects, features, aspects and advantages of the present invention will become more apparent from the following detailed description of the present invention when taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThis invention relates to an attenuator circuit for electrically varying the amount of attenuation of an input signal.\n2. Related Art\nIn general, an attenuator circuit used for a transmitter on the side of the head end (or antenna site) in a two-way community antenna television (CATV) or for a small signal transmitting module such as a communication unit is a .pi. type resistance attenuator which, as shown in FIG. 4(a), is made up of resistors R1, R2 and R3, or a circuit as shown in FIG. 4(b) in which a signal detector circuit 3 detects the output signal of an amplifier 2, and an automatic gain control (AGC) amplifier 1 controls the output gain according to the detection output of the signal detector circuit 3.\nIn the case of the circuit shown in FIG. 4(a), the amount of attenuation is fixed. Hence, the circuit suffers from a problem that the amount of attenuation cannot be changed as desired.\nIn the case of the circuit shown in FIG. 4(b), the amount of attenuation can be changed. However, the circuit is still disadvantageous in that it is high in manufacturing cost because it needs the signal detector circuit made up of a number of circuit elements.\nAccordingly, an object of this invention is to eliminate the above-described difficulties accompanying a conventional attenuator circuit. More specifically, an object of the invention is to provide an attenuator circuit which is able to vary the amount of attenuation, and which is low in manufacturing cost.\nFIG. 7 shows another example of the attenuator circuit of this type, which is made up of a plurality of .pi. type resistance circuits each comprising resistors R1, R2 and R3, and transistor-transistor logic (TTL) drive circuits for controlling the resistance circuits, respectively. For instance, when it is required to attenuate a signal as much as 30 dB in steps of 1 dB, the attenuator circuit should be so designed that the attenuation is made in five steps, 1 dB, 2 dB, 4 dB, 8 dB and 15 dB, as shown in FIG. 7. The attenuator circuit is high in manufacturing cost."}
{"text": "The following discussion, of the background to the invention is included to explain the context of the invention. This is not to be taken as an admission that any of the material referred to was published, known or part of the common general knowledge (in any country) as at the priority date of any of the claims.\nThe banana plant is a large perennial herb with leaf-petiole sheaths that form generally cylindrical, trunk-like pseudostems. Each pseudostem grows from a bud on the true stem (corm), which is an underground rhizome, and can grow to heights normally in the range of 3 to 8 meters over a 9 to 18 month period. When mature, the pseudostem will comprise a soft but dense centre (a ‘core’) surrounded by an outer layer (a ‘sheath’) that is tougher but is less dense, the outer layer typically being formed by the overlapping leaf-petiole sheaths.\nCommercially planted banana plants typically only have 1 to 2 year life-spans, as banana plants only flower (and produce bananas) once, following which the leaves, leaf stalks and pseudostem start to die. This usually requires their removal in some manner, such as by simply being cut down, allowing regrowth of a new pseudostem from the rhizome and the commencement of a new reproductive phase.\nWith annual production in 2002 of about 68 million tonnes of bananas (more than two thirds coming from within India, Brazil, China, Ecuador and the Philippines), it has been recognised that banana pseudostems represent a potentially valuable renewable resource, one which has been traditionally under-utilised and historically economically ignored by banana growers. With this in mind, there have been numerous attempts to use the pseudostems for the production of paper, due to the beneficial properties and qualities of the fibre in the pseudostems.\nHowever, in a paper titled “Banana Stem Fibre for Papermaking” by S. K. Singhai, J. K. Garg and B. Biswas for the Indian Pulp and Paper Journal, August-September 1975, 30(2), pp 13 to 15, the situation at that time was summarised as “The pulping and papermaking qualities of banana (M. sapientum and M. paradisica) stem fibre have been examined. From the available information given in this paper, it is to be considered that neither technically nor economically the use of banana stem fibre is a feasible proposition.” Thus, at least in the mid 70's, no sensible approach had yet been developed to the use of these materials for paper production.\nSince then, it has continued to be recognised that banana pseudostem fibres should have suitable properties for paper production. In a paper titled “Plantain (M. Paradisiaca L) Pseudostem; A Fibre Source For Tropical Countries” by Nicholas A Darkwa of the Forestry Research Institute of Ghana, published in Book 2, TAPPI Proceedings, 1998 for the 1998 Pulping Conference in Montreal Quebec, it was concluded that “ . . . tropical countries that are deficient in long-fibred material for their pulp and paper production can utilise the pseudostems of plantain and banana for such purposes.”\nIndeed, several attempts have been made to use banana plant refuse (predominantly pseudostems, but including leaves, leaf stalks, immature inflorescence and unused bananas) in existing or modified paper pulping processes—see U.S. Pat. No. 5,958,182 for a short summary of some such processes.\nHowever, such refuse commonly has an extremely high water and natural latex content, and includes numerous resinous and gummy substances that are difficult to handle and process. In order to produce workable fibres having desirable characteristics for paper-making, it has proven necessary to extract these fluids and, in particular, wash out the latex and other natural resinous substances. This has proven to be technically difficult, and has generally made the pulping of banana refuse for the production of paper uneconomic, particularly for bulk paper supplies and for anything other than boutique or artistic papers. It has also generally presented the manufacturers with significant chemical waste disposal issues.\nIn Australia, while it has been reported that a good quality paper can be made in low volume by combining banana fibre with that of the betel nut husk (Areca catechu L.), Australian investigators have still concluded that the yield of banana fibre is too low for extraction to be economical. It has been reported that only 1 to 4 oz (28-113 g) of suitable fibre can be obtained from 40 to 80 lbs (18-36 kg) of green pseudostems from the pulping process. Thus, 132 tonnes of green pseudostems would yield only 1 tonne of paper. The conclusion was that the pseudostem would have much greater value as organic matter chopped and left in the field to fertilise subsequent crops.\nThe present invention seeks to provide a method that makes possible the use of banana plants in the family Musaceae for the production of sheets useful in the production of raw paper for subsequent conversion to paper products, ideally in a manner that is both technically simple and reasonably economic, so as to permit relatively high volume paper production therefrom."}
{"text": "Nitric oxide (referred to herein as “NO”) has many uses, including as a medicinal agent. For example, NO has been shown to inhibit smooth muscle proliferation, thrombus formation, platelet aggregation, and smooth muscle contraction. NO can also be used as a bacteriocidal or bacteriostatic agent to sterilize the surfaces of, for example, medical devices. However, the half-life for NO release of most of the known compounds which release NO is less than twelve hours. Thus, most known NO-releasing compounds are too unstable to be useful commercially. The full commercial potential of NO is unlikely to be realized until more stable NO-releasing compounds are developed."}
{"text": "The present invention relates to a method for making a powder of aluminum nitride.\nFor manufacturing an aluminum nitride powder, there are many methods well known to those skilled in the art. For example, Japanese Laid-open Patent Publication No. 53-68700, published in 1978, discloses the manufacture by causing an aluminum compound to react with ammonia or the like to produce an aluminum-nitrogen compound which is in turn heat-treated. Japanese Laid-open Patent Publication No. 54-126697, published in 1979, discloses the addition of an aluminum nitride powder to a powdery mixture of alumina and carbon and the subsequent heat-treatment of the resultant mixture in an N.sub.2 atmosphere. Japanese Laid-open Patent Publication No. 57-22106, published in 1982, discloses heating an aluminum foil and a carbon paper to 600.degree. C. or higher in a non-oxidizing atmosphere, contacting them with an oxygen containing gas, and finally treating them in an atmosphere containing a mixture of oxygen and nitrogen.\nHowever, since all of the prior art methods are of a batch production system utilizing a gas-solid reaction, not only is the productivity low, but also the method is complicated. In the case of the method wherein the aluminum nitride powder is mixed, not only is the running cost high, but also since the powder is used as one of the raw material, the raw material tends to scatter, posing a problem in that the working environment is contaminated."}
{"text": "There is an acute need for new pesticides. For example, insects and mites are developing resistance to the insecticides and acaricides in current use. At least 400 species of arthropods are resistant to one or more insecticides. The development of resistance to some of the older insecticides, such as DDT, the carbamates, and the organophosphates, is well known. But resistance has even developed to some of the newer pyrethroid insecticides and acaricides. Therefore, a need exists for new insecticides and acaricides, and particularly for compounds that have new or atypical modes of action."}
{"text": "Process control refers to the control of the operational parameters of a system by monitoring one or more of its characteristics over time. It is used to insure that the quality and efficiency of the system remain within desired parameters over the course of time. While process control is typically employed in the manufacturing sector for process, repetitive and discrete manufactures, it also has wide application in service industries, such as environmental control.\nProcess control equipment typically utilizes control/sensing devices that are physically integrated into the systems being controlled. For example, a thermostat is typically used in environmental control to insure that building temperatures remain within specified parameters. Likewise, flow control sensors and automated valves are typically used in process manufacturing to insure proper fluid flow volumes.\nThough in early process control systems, control/sensing devices were typically stand-alone units, modern process control systems provide central workstations for monitoring and controlling the control/sensing devices. Particularly robust systems are the I/A Series™ industrial automation systems designed, manufactured and marketed by the assignee hereof, The Foxboro Company, of Foxboro, Mass., USA. In these systems, multiple control/sensing devices are coupled by way of buses to control stations which, in turn, are coupled by way of a local area network (LAN) to one or more operator workstations.\nThe I/A Series systems are built around the client/server model. Client applications software executing on the workstations exchange information with the control/sensing devices via a server, referred to as the “object manager,” executing in distributed fashion in the control stations. Upon request by a client application, the server creates, locates, accesses and updates data structures (“objects”) storing information on the status of at least selected control/sensing devices. For example, a client application that displays temperatures sensed by a thermocouple requests that the server create an object storing a temperature reading from the thermocouple and that the server notify the client each time the temperature changes.\nAlthough modern process control systems, such as the I/A Series systems, have proven quite successful, to date they have provided only limited remote access capabilities. Thus, while numerous operator workstations may reside within the factory or facility in which the control/sensing devices are disposed, it has traditionally proven difficult to access and control those devices outside those areas.\nRemote access and control of processes is desirable for a number of purposes. A plant manager who is “on the road,” for example, may wish to monitor the plant processes while travelling. By way of further example, the manufacturer of process control equipment may require remote access to a plant's control/sensing devices in order to provide technical support.\nAn object of this invention is to provide improved methods and apparatus for process control.\nAnother object of the invention is to provide such methods and apparatus as permit monitoring and control of remote processes.\nStill another object of the invention is to provide such methods and apparatus as can be readily adapted to existing automated process control systems.\nYet still another object of the invention is to provide such methods and apparatus as can be implemented without undue expense and without undue consumption of resources."}
{"text": "The invention relates to a paging receiver comprising a receiving circuit for the receiving and detecting of a first data block transmitted by a transmitting device, of which a second data block represents an address code for one or more receivers;\na first memory in which a received first data block is stored;\na first comparator which detects a first equality if the received address code is equal to an address code assigned to the receiver;\na second memory with two or more memory locations which are each suitable for the storing therein of a third data block of various first data blocks;\na first control circuit which, when the first equality occurs, scans locations of the second memory, compares by means of a second comparator the contents of corresponding sections of third data blocks which are present in the first memory and in the scanned location respectively, and alters the content of the scanned location depending on a second equality detected by the second comparator;\na second control circuit for the selection of a location of the second memory;\nand an indicating element for the indication of a section of the third data block which is present in a location selected by the second control circuit.\nA paging receiver of this type is known from British Pat. No. GB.2,101,779. In the known paging receiver the first, second and third data blocks are equal to each other and equal to the section of the third data block which is read out of the second memory and which can be displayed as a message on a window of the indicating element. On receipt of the data block the first control circuit will scan all the locations of the second memory and compare the content of each scanned location with the information content of the received data block. If an equality is detected in this process, the storage of the received data block in the second memory is prevented. If, after scanning all the locations, no equality is detected the oldest message present in the second memory is erased and the memory space which becomes available as a result of this is used for the storage of the newly received data block.\nIn the known paging receiver a newly received data block cannot be stored in a random location of the second memory. This is in particular troublesome if the oldest message recorded in the second memory contains more important information for the carrier of the paging receiver than the other messages recorded in the first and second memories. The importance of this drawback becomes particularly obvious in paging receiver systems with which emergency calls can be sent to service personnel such as in a nurse paging system.\nAnother important drawback of the known paging receiver is that a message no longer considered of importance by a user of the transmitting device is maintained in the second memory of the receiver so that the carrier of the receiver is inaccurately informed and may take an incorrect action.\nThe object of the invention is to eliminate the drawbacks of the known paging receiver."}
{"text": "In implementing belt or strap type restraining systems, it is sometimes necessary to provide an attachment device to make a connection between an anchoring member and the restraining belts. For instance, the commonly used shopping cart having a child seat (such as might be found at a supermarket or grocery store) may employ a belt restraining system to prevent a child from climbing out of the child seat and causing injury. The restraining system is commonly implemented using belts having opposing buckle pieces attached at an end of each belt. When the buckle pieces are engaged, the belts are fitted tightly around the child's waist, thus restraining the child from exiting the child seat. The buckle pieces can then be disengaged when it is desired to remove the child from the seat.\nIn order to create the restraining effect, the belts used in this system must themselves be anchored to a portion of the shopping cart in order to withstand the applied forces of the restrained object (e.g. a child). It is preferred, however, that the restraining belts of the restraining system not need to be altered in the attachment process, as this adds considerable time and cost to the manufacturing process, and may result in a weak connection. For example, it is uneconomical to stitch the restraining belts to surround an anchoring member of the shopping cart. Also, attaching a restraining belt using a hook and loop connector may result in a weak or easily removed attachment.\nAs a result of these deficiencies, attachment devices are generally used to accomplish the desired anchoring. Attachment devices generally allow the belts or straps of the restraining system to be anchored to the shopping cart by engaging a loop at an end of the restraining strap and by surrounding an anchoring member in the shopping cart. The attachment device creates the connection between the shopping cart and the restraining strap when an interlocking mechanism on the attachment device is engaged, causing the attachment device to completely surround the anchoring member. It is generally desirous that the attachment device provide the strongest possible connection between the shopping cart and the restraining belts or straps in order to prevent failure of the restraining system.\nAttachment devices have been constructed in various forms. For example, an attachment device presently available is constructed from molded plastic and has a V-shaped surrounding member, with two opposing straight bars each attached to either end of the surrounding member. All pieces of the attachment device are solid and cylindrical. The two straight bars terminate with mating portions of an interlock mechanism, such that the straight bars may be locked together to form a connection. When connected, the straight bars form a straight belt member for engaging a loop of a restraining belt, and also thereby surround an anchoring structure within the attachment device.\nAlso available is an attachment device of molded plastic design, formed in approximately a D-shape. A straight bar is used as a belt member to engage a loop at the end of a restraining belt, and a D-shaped surrounding member is coupled at one end to the straight bar such that the attachment device can surround the anchoring structure. A snap interlock is provided on opposite mating ends of the surrounding member and the belt member to lock the attachment device closed once it has been engaged with the restraining strap and the anchoring member of the shopping cart. As with the previous attachment device, the members are also solid.\nDifficulties may arise with these designs which can degrade their effectiveness. The straight bar typically used as a belt member is relatively inflexible, such that when a strong, sharp force is applied by the belt to the attachment device, the attachment device is prone to fracture. Also, cracks can develop in the solid members of the attachment device from other sources (for example from age, weathering, improper use, or manufacturing defects). Whatever the source, cracks can spread through the solid member when force is applied to the attachment device, eroding its resistive properties. Once a crack has extended completely through the solid member, catastrophic failure occurs, and the device must be replaced."}
{"text": "1. Field of the Invention\nThe present invention relates to a power semiconductor device and a manufacturing method thereof, and more particularly to a power semiconductor device having adjustable output capacitance and a manufacturing method thereof.\n2. Description of the Prior Art\nPower semiconductor devices are typical semiconductor devices in power management applications, such as a switching power supply, a power control IC of a computer system or peripherals, a power supply of a backlight, motor controller, etc. Power semiconductor devices can include various kinds of transistors, such as an insulated gate bipolar transistor (IGBT) and a metal-oxide-semiconductor field effect transistor (MOSFET).\nPlease refer to FIG. 1, which is a schematic diagram illustrating a cross-sectional view of a trench MOSFET device according to the prior art. As shown in FIG. 1, the conventional trench MOSFET device 10 of the prior art includes a drain metal layer 12, an N-type substrate disposed on the drain metal layer 12, an N-type epitaxial layer 16 disposed on the N-type substrate 14, two P-type doped base regions 18 disposed in the N-type epitaxial layer 16, two N-type doped source regions disposed in the P-type doped base regions 18, an inter-layer dielectric (ILD) layer 22 and a source metal layer 24. The N-type epitaxial layer 16 has a trench 26, and an insulating layer 28 and a gate conductive layer 30 are sequentially disposed in the trench. The gate conductive layer 30 serves as a gate of the trench MOSFET device 10. In addition, each N-type doped source region 20 and each P-type doped base region are disposed at a side of the trench 26. The N-type doped source region 20 is regarded as a source of the trench MOSFET device 10, and each P-type doped base region 18 adjacent to the insulating layer 28 serves as a channel of the trench MOSFET device 10. The N-type epitaxial layer 16 serves as a drain of the trench MOSFET device 10. The ILD layer 22 is disposed on the gate conductive layer 30 and a part of the N-type doped source region 20, and the source metal layer 24 covers the ILD layer 22, each N-type doped source region 20 and each P-type doped base region 18, so that the source metal layer 24 is electrically connected to the each N-type doped source region 20 and each doped base region 18. The gate conductive layer 30, the ILD layer 22 and the source metal layer 24 form a first gate-source capacitor Cgs1, and the gate conductive layer 30, the insulating layer 28 and each P-type doped base region 18 form a second gate-source capacitor Cgs2. In addition, the gate conductive layer 30, the insulating layer 28 and the N-type epitaxial layer 16 form a gate-drain capacitor Cgd, and a depletion region between the P-type doped base region 18 and the N-type epitaxial layer 16 forms a drain-source capacitor Cds.\nThe desire for ever more compact electronic devices has pushed for size reductions in integrated circuits. Therefore, higher integrations and higher densities are developed continuously. The layout design for the conventional trench MOSFET device 10 has been investigated to reduce the trench width and the trench pitch. However, when the width of the trench 26 is reduced, the coupled area between the gate conductive layer 30 and the N-type epitaxial layer 16 is also reduced, and the contact area of the P-type doped base region 18 and the N-type epitaxial layer 16 is also reduced. Thus a capacitance of the gate-drain capacitor Cgd and a capacitance of the drain-source capacitor Cds are reduced, and an output capacitance of the trench MOSFET device 10 formed by the capacitance of the gate-drain capacitor Cgd and the capacitance of the drain-source capacitor Cds is reduced accordingly. In addition, the capacitance of the drain-source capacitor Cds is far larger than the capacitance of the gate-drain capacitor Cgd.\nGenerally speaking, the conventional trench MOSFET device is usually used in the converter, such as synchronous buck converter, of power management circuit, and is used to be a switching element of the converter, so that the conventional trench MOSFET device usually needs to perform an action of turning on or turning off. When the conventional trench MOSFET device is turned off, the output capacitance of the conventional trench MOSFET device is charged to a same voltage as an outside transformer. However, the converter further includes an inductor device, and the output capacitor and the inductor device form an LC oscillating circuit. Voltage spikes are generated accordingly, and the output capacitance of the conventional trench MOSFET device is reduced with decrease of the total device size and the trench width. Therefore, the voltage spikes caused by turning off the conventional trench MOSFET device are increased, and higher power loss is generated.\nPlease refer to FIG. 2, which is a schematic diagram illustrating a conventional circuit for reducing the voltage spikes. As shown in FIG. 2, a conventional method for reducing the voltage spikes is to electrically connect a snubber circuit 12 in parallel between the source S and the drain D of the conventional trench MOSFET device 10, and the snubber circuit 12 is formed by connecting a capacitor C and a resistor R in series. The capacitor C outside the conventional trench MOSFET device can be used to increase the output capacitance of the conventional trench MOSFET device 10 and reduce the voltage spikes. However, the extra electronic device generates extra circuit cost, and increases extra manufacturing process of welding. Thus the manufacturing cost is increased.\nAccordingly, it is still needed for a novel manufacturing method of power semiconductor device to conveniently and economically resolve or mitigate the problem of the voltage spikes as aforesaid."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nInternal combustion engines draw air into an intake manifold through an induction system that may be regulated by a throttle. The air in the intake manifold is distributed to a plurality of cylinders and combined with fuel to create an air/fuel (NF) mixture. The NF mixture is combusted within the cylinders to drive pistons which rotatably turn a crankshaft and generate drive torque. The drive torque is then transferred to a driveline of a vehicle via a transmission. Exhaust gas produced during combustion may be expelled from the cylinders into an exhaust manifold and then treated by an exhaust treatment system before being released into the atmosphere.\nGases within a cylinder (i.e., the NF mixture and/or the exhaust gas) may enter a crankcase of the cylinder. For example, excessive wear to the cylinder wall and/or the piston ring may allow the gases to enter the crankcase. Gases that enter the crankcase may also be referred to as “blow-by vapors.” The crankcase includes the crankshaft which is connected to the piston. The crankcase also includes oil for lubricating the movement of the crankshaft and the piston. Blow-by vapors may contaminate the oil which may result in damage to the cylinder and/or its components."}
{"text": "The present invention relates to electronic devices, and more particularly to bipolar transistors.\nBipolar transistors are widely used in both digital and analog circuits including computer processors, memories, power supplies, and others. Bipolar transistors are utilized for signal amplification, switching, bias generation, and for other purposes. While bipolar transistors are generally larger and less power efficient than their MOS counterparts, bipolar transistors are typically faster, and consequently bipolar transistors are widely used in fast memories, fast logic arrays, and many other super fast products for data- and telecommunications.\nTypical goals in a bipolar transistor design include a low transistor-generated noise, a high current and power gain (which allow low-power operation), and a high frequency range (which permits a high speed). Another goal is good matching of the electrical characteristics of similar transistors, and in particular good VBE matching measured as the difference between the base-emitter voltages at equal collector currents. VBE matching is particularly important for monolithic circuit design which relies on similar transistors having similar electrical characteristics. These goals are advanced by reducing the transistor base resistance, the emitter resistance, the collector resistance, and the base-collector capacitance. The base resistance, in particular, is a major contributor to transistor-generated noise. Moreover, the base resistance decreases the frequency range of the transistor the transistor unity-power-gain frequency f.sub.max is inversely proportional to the base resistance. See, G. Gonzales, Microwave Transistor Amplifiers (Prentice-Hall, 1984), incorporated herein by reference, page 33. The emitter and collector resistances also contribute to noise and, further, they reduce the current gain and the power gain. In addition, the emitter resistance makes good VBE matching more difficult to obtain, particularly at high transistor currents. Further, the emitter and collector resistances and the base-collector capacitance decrease the transistor frequency range as they reduce both f.sub.max and the unity-current-gain frequency f.sub.T. Thus, reducing the base, emitter and collector resistances and the base-collector capacitance are important goals in transistor design.\nAnother goal is a small size which is needed to obtain a high packing density in the integrated circuit. In addition, reducing the base size decreases the base-collector capacitance.\nFIG. 1 shows in plan view a prior art transistor 110 having three emitters 120-1, 120-2, 120-3 extending across the underlying base 130. The collector (not shown) underlies the base and electrically contacts the collector contact regions 140-1, 140-2. Emitter contact region 150 interconnects the emitters. Base contact region 160 made from doped polysilicon overlying the base surrounds the emitters and extends between the emitters. Contact openings 164-1 through 164-6, 170-1 through 170-8, 180-1 through 180-6 in an insulating layer overlying the transistor allow the emitters, the base and the collector o be contacted by conductive layers overlying the insulating layer.\nThe emitters 120-i are made narrow, i.e., the emitter width WE is made small, in order to reduce the intrinsic base resistance RBI, i.e. the resistance to base current of the base portions underlying the emitters. Of note, base portions underlying the emitters (the \"intrinsic base\") are typically thinner in the vertical cross section and they are typically lighter-doped than the remaining (\"extrinsic\") base portion, and hence the intrinsic base resistance is a significant component of the total base resistance. The number of the emitters--three--is chosen to obtain a desired base-emitter junction area in accordance with the desired emitter current. Emitter contact openings 164-i are positioned away from the emitters not to restrict the emitter width--the contact openings are typically made wider than the emitters to obtain a low contact resistance in the openings. Base contact region portions 160-1, 160-2 extend between emitters 120-i to reduce the base resistance.\nReducing the base, emitter and collector resistances and providing a small size and a low base-collector capacitance are often conflicting goals requiring careful balancing. For example, making the emitters narrow reduces the base resistance but increases the emitter resistance. The base size becomes also increased as more emitters are needed to achieve the desired base-emitter area. Base contact region portions 160-1, 160-2 which extend between the emitters reduce the base resistance but increase the base size. Increased base size leads to a higher base-collector capacitance. The base and collector resistances can be reduced by increasing the base and collector width in the direction along the emitters, but the base size, the base-collector capacitance and the emitter resistance will suffer. Thus, there is a need for a transistor which simultaneously provides low base, emitter and collector resistances, a small base size, a small overall size, and a low base-collector capacitance."}
{"text": "When running electrical or communication lines underground, these lines are often passed through conduits or pipes to help protect the lines from electrical short, fire, explosion, or deterioration or damage due to moisture, rodents, human activity such as digging, shifts in surrounding dirt, gravel or other granular fill and other adverse conditions or activities. In some cases, the pipes may be arranged in an excavation and may further be encased in cast-in-place concrete. That is, an assembly of pipes may be placed in an excavation, formwork may be placed around the assembly of pipes, and cast-in-place concrete may be poured into the formwork to encase the pipes, allowing for electrical or communication lines to be run through the encased pipes.\nThis cast-in-place approach often results in delays in the construction of facilities, due to the need to arrange and wait for the delivery of concrete to the project site. The assembly of pipes may be disturbed between the time of their placement and the concrete pour. If not well-anchored, the pipe assembly may be disturbed by the pour, or poured concrete may enter and occlude the interior of a pipe. In some situations, relatively expensive pump trucks may be required to reach the locations where the cast-in-place concrete is to be placed. Moreover, once the concrete has been placed, the concrete may need to be allowed to cure to reach a specified strength prior to backfilling an excavation, causing further delay. Waiting for concrete test results may also delay the backfilling of an excavation and create additional down time for the contractor."}
{"text": "1. Field of the Invention\nThis invention is directed to an article of manufacture and to a method. More specifically, this invention relates to a foam tile roof panel and system; and, to a method for fabrication thereof.\n2. Description of the Prior Art\nThe use of the synthetic (polymeric) materials as a replacement for more traditional construction materials has and continues to evolve, as the sources of the more traditional materials (wood) increase in costs and will become increasingly scarce.\nAccordingly, plastic or synthetic based building products (once reserved primarily for decorative or cosmetic applications) are, and continue to find uses as structural materials as building codes become more relaxed, or in recognition of the functional equivalency or superiority of the engineering component prepared from the synthetic material.\nStructural foam panels from closed cell expanded foam polystyrene and other common polymers have and continue to be used in lieu of plywood in the construction industry as underlayment for house siding and other construction applications where enhanced energy conservation is desired.\nThe following patents are represented of a use of foam materials (both polymer and concrete) in the fabrication of and in situ manufacture of construction materials: U.S. Pat. No. 3,962,841 (to Carol); U.S. Pat. No. 4,013,614 (to Self); U.S. Pat. No. 4,071,984 (to Larrow); U.S. Pat. No. 4,658,554 (to Riley); U.S. Pat. No. 4,774,794 (to Grieb); U.S. Pat. No. 4,712,349 (to Riley); and U.S. Pat. No. 4,799,982 (to Vicino).\nU.S. Pat. No. 3,962,841 (to Carol) describes an insulating decking structure utilizing a sheet metal structural shape (envelope) containing a perforated foam synthetic board (16) having vertical holes to allow penetration of concrete through the foam board. In brief, the structure of the Carol system comprises a wire reinforced sub-purlin, or purlin, structure that includes a sheet metal envelope (12) or form, containing gypsum foam board (12), a perforated formed synthetic polymer board (13) and poured concrete (15) which penetrates through the polymer board (13) to the gypsum board (12). The exterior surface of the foregoing composite is finished with alternating layers of roofing felt and tar (16) and a waterproof wearing surface of tar and gravel (17).\nU.S. Pat. No. 4,013,614 (to Self) describes a method for preparation of spray-up laminates by foaming-in-place an unsaturated polyester resin (hereinafter \"synthetic resin\"). The synthetic resin can be applied on both sides of expanded metal lathe (Example 3) and applied directly to corrugated metal building panels having a spray applied polyurethane base (example 4). The fabrication of laminated panels (in accordance with the teachings of Example 4) further contemplate a sequential application of these materials to the structure of metal panel (column 5, lines 36-58).\nU.S. Pat. No. 4,071,984 (to Larrow) describes a structure, such as a house, assembled from pre-fabricated panels. The wall panels of the structure comprise a laminate, or composite, structure having an interior section of foamed material, such as polystyrene or urethane foam, with another surface comprising wood, aluminum or plastic. The apparent invention resides in the configuration of the panels to allow for their interlocking or integration with one another in the construction of a shelter.\nU.S. Pat. No. 4,658,554 (to Riley) describes an insulated roofing system in which a waterproof plastic film or membrane is bonded to the under surface of a panel to positively exclude water. The panels to which this membrane is applied are preferably foam can be provided on a its upper surface with integral ribs, such as grooves, and of sufficient structural integrity to support concrete panels which can be placed upon and supported by the integral ribs. The integral ribs which are positioned between the foam panel and the concrete surface member provide (what is referred to in the Riley patent) \"channels\" for ventilation and moisture removal on hot, dry days. The concrete panels are not, according to Riley, are \"simply abutted against one another, and preferably are laid so that the butted junctions do not coincide with the butted junctions of the underlying foam panels\" (column 3, lines 31-34). Thus, it is apparent from Riley that the waterproof membrane on the interior of the composite is necessary to exclude water which penetrates between the concrete block and is not otherwise carried away by the ribs (16) which are intermediate between the concrete and the styrofoam panels. To the extent that any moisture does intrude beyond the interior ribs, the waterproof membrane (12), positioned between the ribs and the foam will presumably prevent any further intrusion thereof.\nU.S. Pat. No. 4,774,794 (to Grieb) describes a foam cement building having walls, roof and/or floor formed from a plurality of self-supporting foam building blocks of varying densities. The exterior surface from the foam building blocks are coated with an architectural finish-such finish being formed from cement, reinforced with fiberglass mesh or fiberglass rovings or strands. The blocks are connected to one another by a \"tongue and groove\" arrangement which provides, once assembled, a monolithic structure.\nU.S. Pat. No. 4,712,349 (to Riley) is essentially the same, (at least in relation to the instant invention), as U.S. Pat. No. 4,658,544 (discussed above).\nU.S. Pat. No. 4,799,982 (to Vicino) describes a method and system for the manufacture of a modular, monolithic molded building structure from a dismantleable mold. The Vicino patent describes the sequential application of an expandable foam to various mold components along with reinforcing materials. The purpose and effect of the process described by Vicino is to provide a seamless building structure, (with the exception of windows and doors).\nAs is evident from the foregoing discussion, the techniques utilized in the past for the fabrication of foam based building products, including those prepared from synthetic polymers and concrete, has involved a variety of different expedients to accomplish the specific objective of their respective inventors. As is further apparent, none of the art discussed above provides a building product having both structural and finished features incorporated into a single panel, or panel array, suitable for use as a roofing material. In virtually all of the structures described above in the prior art, the use of an expandable foam has been simply to provide an insulating element without appreciation of its versatility; specifically the combination of its aesthetic and structural properties in a single element, while at the same time eliminating other unnecessary components which can reduce costs, construction complexities and weight, (which can be a material factor in engineering the trusses and other load bearing components)."}
{"text": "Recently, a large-capacity magnetoresistive random access memory (MRAM) using a magnetic tunnel junction (MTJ) element has been expected and attracted attention. In the MTJ element used for the MRAM, one of the two ferromagnetic layers that sandwich a tunnel barrier layer is set as a magnetization fixed layer (reference layer) where the direction of magnetization is fixed, and the other layer is set as a magnetization free layer (storage layer) where the direction of magnetization is made to be easily inverted. In this structure, information can be stored by associating, with binary numbers “0” and “1,” a state in which the direction of magnetization is parallel between the reference layer and the storage layer, and a state in which the direction of magnetization is opposite between them."}
{"text": "1. Field of the Invention\nThe present invention relates to a battery remaining amount warning circuit applicable to an electrical product or electronic device using a battery.\n2. Description of the Related Art\nToday, many electrical products and electronic devices use batteries. As a conventional method for warning about a remaining amount of a battery, it is a general practice that a battery voltage is measured and a warning about the remaining amount is issued when the battery voltage falls below a specified value or DC resistance of the battery increases.\nThere are a variety of types of battery such as alkaline cell, nickel hydrogen cell, nickel cadmium cell, and there are also a variety of battery characteristics. However, with a monitoring of the battery voltage alone, depending on the type of the battery, the battery power may not be used to the full because a considerable amount of the battery can still be used after the warning is given about the insufficient remaining amount of the battery.\nFurthermore, there is also a proposal of a method for constantly measuring a battery voltage by providing a dummy-load circuit (Japanese Patent Publication No. 63-3538). However, this method increases unnecessary power consumption.\nThere is also a proposal of a method for deciding whether a camera operates normally or not based on a boosting operation of a booster circuit when a flash capacitor is charged (Japanese Patent Application Publication No. 2000-56370), but when no flash is used or can be used, there is no way to decide it.\nFurthermore, when the remaining amount of the battery is predicted from the battery voltage, the voltage varies depending on the type of the battery and when the battery voltage is low because of the nature of the type of the battery, a warning is issued earlier. For example, a warning is issued with an alkaline cell at 1.5 V/cell and with a NiCd cell at 1.2 V/cell."}
{"text": "The information provided below is not admitted to be prior art to the present invention, but is provided solely to assist the understanding of the reader. All U.S. patents and patent applications and other published documents mentioned anywhere in this application are incorporated by reference in their entirety.\nThe family of semi-automatic firearms based on the original ARMALITE AR-15 and COLT M16 rifles has achieved significant success with military, law enforcement and civilian shooters alike and is currently one of the broadest and most popular in the world. improvements to the AR-15/M16 platform are diverse in their purpose, with new embodiments and aftermarket parts and accessories continually being developed due to the overwhelming popularity and inherent modularity of the base platform. Although the template firearm design itself has undergone many refinements to improve its functional reliability, to expand its mission roles, and to reduce its manufacturing costs, the design's principles of operation have not been significantly improved. In fact, firearms patterned after the reputable AR-15/M16 template using the same basic design features and specifications are currently produced in numerous calibers and countless versions and configurations by different manufacturers throughout the world.\nThe term “AR-15®” as used herein refers to the rifle manufactured by COLT DEFENSE LLC or its licensees at the time of invention. Semi-automatic variants and derivatives of the AR-15/M16 produced by other manufacturers are usually termed “AR-15 type” or “AR-15 style” regardless of caliber and are sold under various proprietary model designations. The term “AR type” is used in this disclosure to refer to authentic semi-automatic COLT firearms based on the AR-15/M16 platform such as the AR-15, AR-15 Sporter SP1, AR-15A2 Match Target and AR6720, as well as semi-automatic versions of AR-15 type firearms such as the ARMALITE AR-10, REMINGTON R-15 and the BUSHMASTER XM-115 rifles at the time of invention.\nThe use of lighter weight synthetic materials such as reinforced composites and polymers (collectively, “polymers” or “synthetics”) to form various parts of a firearm such as a receiver or frame for a handgun or rifle has increased significantly in recent years. Lighter weight synthetic materials provide the firearm with a reduced weight for ease of carrying, handling and use in the field, on the range and during combat. For example, by reducing the weight of a firearm, the user generally is able to manipulate the firearm faster and easier, such as while tracking moving targets. The lighter weight also means less load that must be borne by the user during transport and use, which conserves the user's energy, reduces fatigue and increases accuracy. Additionally, most lighter weight synthetic materials now being used for firearms are typically less susceptible to corrosion and damage resulting from exposure to water, salt, dirt, and other environmental hazards. Moreover, many modern polymers are easier and less costly to manufacture into firearm components than traditional metal substrates such as forged aluminum, aluminum billet or steel, making components formed from synthetic materials more affordable.\nA major problem associated with the use of lighter weight polymer materials to form parts of a firearm, particularly receivers and actions, is that such materials sacrifice strength and the ability to withstand the extreme pressures and stresses created upon firing a round of ammunition in exchange for lighter weight. For example, upon discharge of an AR type firearm, the basic design of which is well known in the art, the bolt is forced rearward to an open position to eject the spent cartridge and ready the chamber to receive a new round, after which the new round is loaded in the chamber by the capture and backflow of gases created upon the firing of the previous round of ammunition. The extreme chamber pressures and recoil forces resulting from the travel of the bolt are primarily borne by the receiver of the firearm. Over time, such extreme stresses can cause cracking and potentially failure in a receiver formed from some lighter weight synthetic materials. For example, the primary point of stress and the most likely point of failure in an all polymer lower receiver for an AR type firearm is around the threaded ring or collar disposed at the rear end of the lower receiver body that is configured to accept and connectively engage with a threaded receiver extension (also commonly and interchangeably known as a “buffer tube”). Much of the stress caused by the recoil force of a discharging AR type firearm is focused around the point where these two components engage.\nA further related consequence of using lighter weight, polymer-based materials such as plastics and reinforced composites to form synthetic firearm receiver bodies has been that such materials typically are not as durable as conventional forged steel or aluminum, the traditionally preferred substrates for firearm receivers. More pointedly, the threaded portions of all-polymer receivers tend to wear out more quickly, especially when connectively engaged by conventional metal screws or screw threads, which leads to a higher chance of failure and thus an increased risk of firearm malfunction and shooter injury during use. This defect is of particular concern to users of AR type polymer lower receiver bodies (as used herein, the term “lower receiver body” means a lower receiver stripped of all its operational components, also commonly known in the art as a “stripped” lower receiver) because many of the currently available parts, components and accessories for AR type firearms, such as the buffer tube and pistol grip retaining screw, are still made from one or more types of metal and are configured to engage the lower receiver body using screw threads. No suitably strong, yet lightweight AR type stripped lower receiver durable enough to lastingly withstand recoil stress and engagement with metal screws is currently available.\nMore durable, lighter weight AR type stripped receivers, including upper and lower receivers and lower receiver assembly kits are still being sought. As a result, it is desirable to provide a threaded metal insert for a polymer lower receiver body, particularly for an AR type polymer lower receiver such as that of an AR-15 or AR-10, that provides the greater durability and increased strength and longevity of traditional metal alloys but which does not substantially affect the lighter weight imparted by the polymer. It is also desirable to provide an inexpensive, hybrid polymer and metal stripped lower receiver for an AR type firearm such as an AR-15 or AR-10 that overcomes the foregoing shortcomings of the prior art to deliver a lighter weight, corrosion resistant, stripped lower receiver having superior strength and durability that functions with substantially all high quality conventional AR type (upper and lower) receiver parts kits, yet does not require the use of any specific caliber, brand or type of parts. It is also desirable to provide methods for economically manufacturing such threaded metal inserts and hybrid polymer lower receivers. Accordingly, a need exists for an improved polymer lower receiver for AR type firearms that addresses the foregoing and other related problems in the art to combine the lighter weight and corrosion resistance of all-polymer lowers with the greater strength and durability of traditional metal components. Other objects and advantages will become apparent from the following disclosure."}
{"text": "The background art discussed in this section includes information disclosed under 37 CFR §1.97 and §1.98.\nThe performance of current directional measurement instruments, such as network analyzers, is inadequate for many possible applications that extend significantly above one hundred gigahertz. Existing and proposed directional measurement systems that form the cores of directional measurement instruments suffer from limited directivity or bandwidth at frequencies above one hundred gigahertz.\nA lightwave is modulated by a modulation wave in an optical modulator. The modulation wave is often in the form of a radio wave, microwave, or millimeter wave and guided by a modulation waveguide of the optical modulator. A travelling-wave optical modulator may be employed in an optical system having modulation in the microwave or millimeter-wave frequency range. In a travelling-wave optical modulator, the modulation wave is substantially coherent with the lightwave modulation over the length of the modulator, and operation is said to be phase-matched."}
{"text": "High speed, high density, and high endurance memory device are desired in a modern computer system in order to improve the overall performance of the computer system. Dynamic random-access memory (DRAM) has low latency, short access time, and high endurance. However, the size of DRAM in a computer system is generally limited due to, for example, relatively high cost and difficulty in making large DRAM chips or modules. An electrically accessible non-volatile or persistent memory, often referred to as storage-class memory (SCM), that combines the benefits of DRAM, such as high performance and robustness, with the persistent capabilities and low cost of conventional solid-state storage has been developed and is beginning to be used in modern computer systems."}
{"text": "1. Field of the Invention\nThe present invention relates generally to packaging, and particularly to a packaging arrangement for securely containing one or more light-sensitive articles in a light-tight environment.\n2. Description of the Prior Art\nPackaging arrangements for containing various types of articles are well known in the prior art. Examples may be found in the following patents:\nU.S. Pat. No. 2,942,768 (McCall)--discloses a rectanguloid box made from a single blank of cardboard folded to form a box body and a foldable lid, the lid comprising a top-wall section whose opposite ends overhang the body end walls and a narrower, foldably depending tucker section that closely faces the inside surface of the body front wall and whose opposite ends lie close to the end walls, each end wall comprising outer, intermediate, and inner end-wall panels, the outer end-wall panel having an inwardly folded extension thereof that passes beneath the intermediate and inner end-wall panels toward the body interior and overlies the bottom wall.\nU.S Pat. No. 3,269,638 (Forbes)--discloses a rectanguloid box made from a single blank of cardboard folded to form a box body and a foldable lid, the lid comprising a top-wall panel whose opposite ends overhang the body end walls and a narrower, foldably depending tucker panel that closely faces the inside surface of the body front wall and whose opposite ends lie close to the inside surfaces of the body end walls, each end wall comprising outer, intermediate, and inner end-wall panels, each outer end-wall panel having an inwardly folded bottom panel that extends therefrom beneath the adjacent intermediate and inner end-wall panels, toward the body center, and overlies the body bottom wall.\nU.S. Pat. No. 4,295,599 (Locatelli et al.)--discloses a rectanguloid box made from a single blank of cardboard folded to form a box body and a foldable lid, the lid comprising a top cover panel, whose opposite ends overhang opposing body end walls, and a narrower, beveled, foldably depending flap panel that frictionally bears against the inside surface of the body front wall, and whose opposite ends frictionally bear against the inside surfaces of the body end walls, to effect a seal between the cover panel and the front wall and to keep the cover panel closed, each of the body end walls comprising an inverted-U-shaped continuum of outer, upper, and inner end-wall panels extending from and secured to the adjacent end of the body bottom wall.\nU.S. Pat. No. 4,027,794 (Olson)--discloses a rectanguloid container for shipping a cylindrical object having an axially disposed core, the container comprising a rectanguloid inner box including complemental base and cover sections that enclose the object and a mating outer box that is bottomless and telescopes downward over the inner box.\nWhile prior-art arrangements such as those described above may have sufficed for their own particular purposes, there has remained, nonetheless, a need for an improved packaging arrangement that is especially adapted to light-tightly secure a light-sensitive article, such as a photosensitive web roll, especially such a roll that is of relatively large size, without requiring a separate opaque wrapper around the light-sensitive article itself.\nAn overall object of this invention, therefore, has been to provide a packaging arrangement that meets the foregoing need, and to do so in an efficient, reliable, and convenient manner."}
{"text": "Photography and moviemaking have become staples in modem society. A photograph is typically a single image frame of a real scene. Movies can be described as a series of frames that together form what appears to the human eye to be a continuously moving image. Both photographs and movies are now found in both physical and digital formats.\nImage processing software is currently used by millions of people worldwide to perform a variety of image processing functions, perhaps the most typical of which is manually changing colors of pixels to make the image appear more realistic. Examples include removal of “red-eye” from human subjects, and touching up of edges.\nHowever, image processing software has heretofore been unable to accurately identify shadows in an image. Nor has it been able to accurately tie shadowed portions of an object to nonshadowed portions of the same object. As such, the ability to work with shadows in image processing software is limited and inaccurate at best.\nWhat is therefore needed is a way to automatically and accurately identify shadows in an image, and tie shadowed portions of objects to nonshadowed portions. This would open the door for improved photo processing capabilities, and allow automation of features that now require the manual user input. Such a solution would also reduce the inherent flaws in the effects heretofore manually created by human users.\nIt would also be desirable to accurately remove, lighten, or darken shadows from images. Prior methods of removing or lightening shadows merely lighten the entire image, giving the image a “washed out” look. Similarly, prior methods of darkening shadows darken the entire image. These methods fail to maintain the integrity of the nonshadowed portion of the image. Accordingly, the user is often required to revert to manual tools to adjust any shadowing."}
{"text": "1. Field of the Invention\nThe invention is directed to a clutch release arrangement, particularly for a clutch in a motor vehicle, including a master cylinder for a hydraulic system, a pedal for actuating the master cylinder to disengage the clutch, and an articulated energy accumulator which stores energy when the clutch is engaged and provides a boosting force which reduces the actuating force when the clutch is disengaged.\n2. Description of the Related Art\nRelease arrangements of this type are frequently used in modern hydraulically actuated clutches and are known from practice. The energy accumulator has a spring which holds the pedal in its basic position or at least reinforces a force holding the pedal. When the pedal is depressed, the spring is pretensioned slightly. In so doing, the spring swivels until its point of application of force at the pedal goes beyond a dead center and reinforces the movements of the pedal by an added force on the pedal path following the dead center. This enables a reduction in the force expended when depressing the pedal, at least along the path following the top dead center, where the energy accumulator adds its energy as a boosting force to the pedal force.\nDue to the articulation of the energy accumulator at the pedal, in particular as an over-center spring, the pedal force on the second half of the pedal path can be noticeably reduced in order to be able to overcome more easily the spring acting in the clutch of the motor vehicle, particularly at the end of the pedal path. However, known kinematic couplings of energy accumulators with a clutch pedal have shown themselves to be in need of improvement for certain applications. Above all, it is disadvantageous that the path of the force reinforcing the pedal force can not only be adjusted in height with sufficient accuracy, the variability of the reinforcement force along the pedal path can also not be adapted to the requirements of the pedal force on the pedal path. For example, when the maximum force should be further reduced, this can only be done in the known solution by increasing the force for depressing the pedal."}
{"text": "In the manufacture of liquid crystal panel elements and semiconductor elements such as IC and LSI, a chemically vapor-deposited (CVD) electroconductive metal film and an insulating film such as SiO2 film are firstly prepared on a substrate such as silicon wafer or glass. After that, a photoresist is uniformly applied on the electroconductive metal film and the insulating film and then selectively exposed to light and developed to form a resist pattern. Then the substrate is selectively subjected to a dry etching process using the resist pattern as a mask whereupon a fine circuit is formed. After that, photoresist film residue and etching residue which are no longer necessary are washed and removed using a stripping liquid.\nAlthough some stripping effects are achieved when an aqueous solution of sodium hydroxide or a common organic solvent is used solely as the stripping liquid, the resulting stripping effect is not sufficient. Accordingly, various stripping liquids for the photoresist have been proposed already for enhancing the stripping property. One of the common methods among them is a method where a solution of quaternary ammonium hydroxide such as tetramethylammonium hydroxide (TMAH) is used.\nIn that method, preparation of a stripping liquid is usually carried out in such a manner that a solution prepared from commercially available water-containing crystals of quaternary ammonium hydroxide or a commercially available aqueous solution of quaternary ammonium hydroxide is diluted with an organic solvent such as dimethyl sulfoxide (DMSO) or 3-methoxy-3-methyl-1-butanol (MMB) down to a desired concentration.\nOn the other hand, it has been known that the stripping property of a stripping liquid is also dependent upon the amount of water contained in the stripping liquid and that the lower the water content, the higher the stripping property. Since a stripping liquid is prepared by diluting a quaternary ammonium hydroxide solution with an organic solvent as mentioned above, it is necessary for lowering the water content in the stripping liquid that a quaternary ammonium hydroxide solution is concentrated so as to make the water content therein low.\nWith regard to a method for concentrating a quaternary ammonium hydroxide solution to make the water content low, there is a known method whereby a pentahydrate of TMAH (ratio by weight of TMAH to water is 1:1) is dissolved in methanol and concentrated in vacuo according to a batch system and also a method whereby a DMSO solution containing TMAH pentahydrate is dried by adding molecular sieves thereto (please see Patent Documents 1 and 2).\nHowever, in the former method, there are disadvantages that heating for long time is necessary whereby TMAH is apt to be decomposed and that, since the product obtained after the concentration is a solid, its handling is difficult. On the other hand, in the latter method, there are disadvantages that, due to the use of molecular sieves, the cost is high and the operation is complicated. Further, since the presence of fine particles in the concentrated solution and in the stripping liquid is not preferred, there is a disadvantage that the molecular sieves must be removed after drying. Furthermore, the water content in the concentrated solution prepared by such a method is not satisfactorily low. Accordingly, the conventional process for producing a concentrated solution of quaternary ammonium hydroxide still needs to be improved."}
{"text": "Field\nThe disclosed technology relates generally to semiconductor devices and more particularly to three dimensional semiconductor memory devices, such as vertical three dimensional non-volatile memory devices.\nDescription of the Related Technology\nRecently, ultra-high density storage devices have been proposed. Some of the proposed devices are formed using a process architecture comprising a vertical channel three-dimensional (3D) stacked memory structure. One of the process architectures is referred to in the relevant industry as Bit Cost Scalable (BiCS) architecture. Some 3D NAND stacked memory devices are formed using a process architecture comprising a stack of alternating conductive and dielectric layers.\nHowever, various process architectures for forming 3D NAND stacked memory devices take into account a fundamental trade-off: in order to maximize the density, one should share as many gates of memory cells in a given word plane as possible. In this regard, some process architectures use a stack of alternating polysilicon layers and dielectric layers, which are sometimes referred to as a polysilicon-oxide-polysilicon-oxide (POPO) approach, where P stands for polysilicon and O stands for silicon oxide (e.g., SiO2). However, one of the main drawbacks of the POPO approach is the relatively high resistance of the word lines formed by the alternating polysilicon layers. Alternatively, some process architectures contemplate using a metal gate electrode instead of a polysilicon gate electrode. This approach is sometimes referred to as a metal-oxide-metal-oxide (MOMO) approach, where M stands for metal and O for silicon oxide (e.g., SiO2). However, the MOMO approach has been hampered because of the relatively high stress in the stack.\nAnother approach is to apply replacement gate processes in which a sacrificial gate layer, e.g., a layer of silicon nitride, is used to form the alternating stack. Afterwards this sacrificial gate layer is replaced with the desired gate electrode material. This so-called nitride-oxide-nitride-oxide (NONO) approach however, where N stands for silicon nitride (e.g., Si3N4) and O for silicon oxide (e.g., SiO2), sometimes limits the width of the word plane to a few, e.g., four, cell strings. Such replacement gate processes have limited capability in filling the cavities of the removed replacement gate layer. This limited filling capacity can decrease the memory density because slits, through which the desired gate electrode material is provided, are formed every few rows. In addition, these slits have to go all the way down to the substrate which results in a large area. US2015/0079742 (see FIG. 6) discloses such a vertical channel-type three-dimensional semiconductor devices having replacement gates.\nThus, there is a need for novel and improved three-dimensional non-volatile memory devices, and fabrication methods thereof, with a reduced resistance-capacitor time constant of the word line, but still having a high memory density."}
{"text": "1. Field of the Invention\nThis invention concerns a device for assistance with interception by an aircraft for a flight path segment located in a horizontal plane, and a system for assistance with the interception and follow-up of such a segment.\n2. Description of Related Art\nIn general, when an aircraft pilot, for instance flying a helicopter, wishes to intercept a flight path in the horizontal plane, he has to perform the following operations:                identify the flight path part (segment) to be reached, and take a heading for its interception allowing for the effect of wind, if known;        anticipate the turning of the aircraft to join the flight path without overshoot, depending on the speed of the aircraft and the angle of interception; and        ensure that the interception flight path is reliable with respect to ground and any nearby obstacles.        \nThe precision required to follow an intercepted flight path, including a multitude of segments, varies along said flight path. Accordingly, the flight path monitoring accuracy is far greater during final approach than during cruise flight phases.\nIt is known that the pilot in the approach phase in a horizontal plane uses ILS (Indicator Landing System) type information, displayed by dedicated symbology of varying sensitivity, which is constantly increasing as the aircraft approaches the touchdown point. Conversely, during the cruise phase, in order to provide sufficient accuracy, a flight path segment monitoring procedure in the horizontal plane requires the pilot to select the most appropriate scale at a given moment, so as to comply with the required precision and/or use flight path deviation digital information combined with another dedicated symbology, for instance an HIS (Horizontal Situation Indicator).\nAccordingly, the pilot following the current flight phase must make various selections and gather appropriate flight path deviation information to follow at various positions of the display screen (and even on two different screen pages), which information is generally presented using different formalisms.\nFurthermore, during the periods of intense workload encountered by the pilot, omitting to modify a scale can have damaging consequences. Indeed, the pilot may believe erroneously that the aircraft is properly aligned on the flight path to follow whereas, in reality, the selected scale is too large.\nFinally, in order to accomplish with satisfactory accuracy, the complete procedure demands a great deal of know-how and also intense attention and activity (high workload) throughout the interception and follow-up duration. Note that organization and dosing difficulties are essentially due to the fact that the customary instruments and/or symbologies used are not organized to present an obvious synthesis of the effective situation.\nThis invention concerns a device assisting with interception by an aircraft of a flight path segment located in a horizontal plane designed to remedy these drawbacks.\nIndeed, according to the invention, said device assisting with interception by an aircraft of a flight path segment located in a horizontal plane, said device being brought abroad the aircraft and including:                initial resources for determining the values of the aircraft parameters; and        initial display resources to present on at least one display screen:                    at least said flight path segment in the horizontal plane to be intercepted; and            an initial symbol illustrating the position of the aircraft in the horizontal plane with respect to said flight path,is outstanding in that the initial resources determine at least one ground speed vector representative of the speed of the aircraft with respect to the ground and in that said initial display resources are liable to present, in addition, on said display screen:                        a first means of indicating the ground speed corresponding to a first rectilinear section in said horizontal plane connected to said symbol whose angular orientation corresponds to the direction of said ground speed vector and whose length:                    is proportional to the ground speed vector module when said ground speed is greater than a predetermined value; and            is fixed and proportional to the module of said predetermined value when said ground speed is less than or equal to this predetermined value;                        a second means of interception approach indication, corresponding to a second rectilinear in said horizontal plane, in the direction of said flight path segment to be intercepted and whose length adjusts automatically so that said second rectilinear section extends at one of its ends, said first rectilinear section and at the other of its ends, a third means of indication; and        said third means of interception turn indication, corresponding to a curve portion including at least one circle arc in said horizontal plane, one end of said curve portion extending at a tangent into second rectilinear section and, at the other end, connected at a tangent to said flight path to be intercepted.        \nAll this graphic information forms the flight path interception path which ends, if there is no flight path overshoot, in an interception turn (circle arc more or less deformed to allow for wind effect) in the horizontal plane. This graphic representation enables the pilot to see and act early enough to adjust his approach, then his interception turn. This results in a gain in terms of efficiency and availability because, in particular, it relieves the pilot of the need to release the controls.\nIn addition, thanks to the first means of indication, the pilot is given graphic information supplying him both with the real direction toward which the aircraft is moving and a first estimation of the ground speed vector module at which it is moving. In addition, this indication offers the advantage of being displayed close to a screen position used by the pilot a great deal during the pre-defined flight path interception/monitoring phase, i.e. said first symbol indicating the current position of the aircraft.\nIn an advantageous embodiment, said first means of display present the second means of indication on said display screen only:                if a segment of the flight path is in the direction of said first rectilinear section of first indication resources; and        if this segment is at a distance less than a predetermined distance with respect to said first rectilinear section,        \nThanks to said second means of indication, the pilot has synthetic graphic information giving him the future rectilinear interception flight path if he continues to fly with the same current ground speed vector (module and direction). As soon as the ground speed vector varies, said graphic information is updated in real time.\nIn addition, in an advantageous embodiment, said circle arc of the third means of indication includes a radius depending on the speed of the aircraft and/or the aircraft capacities. Furthermore, said circle arc may be deformed to allow for the intensity and direction of the wind.\nFurthermore, in an advantageous embodiment,                if a flight path segment is in the direction of said second rectilinear section of the second indication resource;        if this segment is located at a distance less than a predetermined distance from said second rectilinear section; and        if the flight plan requirements make an interception maneuver possible.        \nThanks to the third indication resource, the pilot has synthetic graphic information supplying him with the future turn to be made to intercept with and align on the flight path if the aircraft continues to fly along the same ground speed vector (module and direction).\nIn one specific embodiment, said interception assistance device includes, in addition, second means for detecting an obstacle and the first means of display are liable to present, in addition, on said display screen, at least a second symbol illustrating the position of said obstacle in said horizontal plane.\nThis invention also concerns an on-board system for assistance with the interception and monitoring by an aircraft of a flight path segment located in a horizontal plane. According to the invention, said system is outstanding in that it includes:                an interception aid device as specified above; and        a flight path follow-up aid device.        \nIn addition, according to the invention, said flight path monitoring assistance device includes:                third means of determining:                    a lateral deviation of the aircraft from a segment of the flight path in the horizontal plane; and            lateral margins tolerated either side of said flight path segment in the horizontal plane; and                        second means of display liable to appear on the display screen:                    a fixed (immobile) lateral deviation scale;            a fourth means of effective lateral deviation indication corresponding to a straight line segment provided for on said scale and illustrating the effective lateral deviation of the aircraft from the flight path segment to be followed; and            a fifth means of indicating excessive lateral deviation appearing on said scale when the aircraft approaches at a predetermined distance, one of said lateral margins.                        \nIn an advantageous embodiment, the size of the fixed scale remains constant thanks to the automatic adjustment of the proportion between the real deviation and its representation on the display screen. Accordingly, whatever the aircraft flight phase, the fixed scale preserves the same size on the display screen.\nIn addition, in an advantageous embodiment:                said straight line segment of said fourth means of indication is inclinable; extraordinary; and/or        when said fourth means of indication reaches its stop, it changes form and remains at the stop; and/or        said fifth means of indication is flashing.        \nTo assist with the interception of a flight path including a multitude of segments, advantageously, said second means of display will also present on said display screen a sixth means of indicating a change in flight path segment to follow.\nIn addition, advantageously, said second means of display will present on said display screen a seventh means of indicating the lateral velocity tendency.\nIn a simplified and preferred embodiment, said first display resources for the device providing interception assistance and said second means of display of the second monitoring assistance device will be part of one and the same display unit.\nFurthermore, the system in accordance with the invention will also include means of transmitting information liable to connect said first means of said interception aid device and said third means of said monitoring aid device to an acceptance autopilot for the transmission of information. Accordingly, all the data generated by the first and third means can be transmitted to an autopilot, which deals with automatic aircraft guidance."}
{"text": "1. Field of the Invention\nThe present invention relates to a vaporizing and supplying method intended for supplying a semiconductor manufacturing apparatus with a gaseous CVD (chemical vapor deposition) material. More particularly, the present invention pertains to a vaporizing and supplying method which is intended for supplying a semiconductor manufacturing apparatus with a liquid CVD material by controlling the flow rate of a liquid CVD material with an extremely high accuracy and vaporizing the liquid CVD material, and is also intended for enabling the production of a semiconductor thin film with an extremely high accuracy of film thickness.\n2. Description of the Related Arts\nAlong with the development of semiconductor industries in recent years, there has been much progress in enhancement of performance and integration for a semiconductor device. Thus, a variety of organometal compounds in liquid form come to be used as materials for metallic films and insulation films to take the place of gaseous hydrides and halides that have heretofore been employed.\nFor instance, there are employed in metallic films for a semiconductor device, dimethyl aluminum hydride {(Al(CH3)2H} as CVD materials of aluminum films; hexafluoroacetylacetone copper vinyl trimethylsilane {(CF3CO)2CHCu.CH2CHSi(CH3)3} as CVD materials of copper films; bis(ethylcyclopentadienyl)ruthenium {(Ru(C5H4C2H5)2} as CVD materials of ruthenium films; etc.\nWith regard to an insulated film for a semiconductor device, there are known SiO2 as a gate insulating film, Si3N4 as a capacitor insulating film and PSG (phosphorus/silicon/glass) and BPSG (boron/phosphorus/silicon/glass) as an interlaminar insulating film. In addition, there are used tetraethoxysilicon {Si(OC2H5)4} as a CVD material for SiO2 film, trimethoxyboron {B(OCH3)3}, trimethoxyphosphorus {P(OCH3)3} and the like as a CVD materials for PSG and BPSG films.\nOn the other hand, as a method for vaporizing a liquid CVD material and supplying a semiconductor manufacturing apparatus with the resultant gaseous CVD material, there has been used a method in which a carrier gas is introduced in a liquid CVD material so as to vaporize part of the liquid CVD material, and the mixed gas of the carrier gas and the CVD material is subjected to flow rate control with a mass flow controller and is supplied to a semiconductor manufacturing apparatus. However, this method is difficult to stabilize vaporization conditions and produce a feed gas at a definite concentration and a flow rate. Accordingly, there is prevailingly employed at the present time, a method in which a liquid CVD material is subjected to flow rate control with a liquid flow controller, and supplied to a vaporizer, where it is vaporized, and the resultant gaseous CVD material is supplied to a semiconductor manufacturing apparatus.\nAs a method for subjecting a liquid CVD material to flow rate control, there are available a method in which an inert gas is pressurized into a liquid CVD material vessel, whereby the liquid CVD material is introduced into a liquid mass flow controller, and is subjected to flow rate control; and a method in which the flow rate of a liquid CVD material is controlled with a flow variable micropump. In any of the above-mentioned methods, a CVD material is supplied to a vaporizer after the flow rate thereof is measured in the state of liquid. Among the methods, there is generally carried out a vaporizing and supplying method using a liquid mass flow controller which is relatively easy in maintenance work.\nA mass flow controller for general purpose used for the flow rate control of a liquid CVD material is constituted of a flow rate sensor, an electrical circuit, a control valve and the like. The flow rate sensor is of such a constitution that two self-heating type resistors wound on the outside of CVD material piping are incorporated in a bridge circuit, and is in such a state that the resistors are heated by a current flowing at all times. The bridge circuit is constituted so that voltage is generated when the balance of the self-heating type resistors is lost. The voltage of the bridge circuit which is caused by the heat transfer due to the flowing of a CVD material in piping is directly proportional to the mass flow rate of the CVD material, and accordingly the flow rate control of the CVD material with the mass flow controller is carried out, for instance, by measuring the flow rate of the CVD material through measurement of electric output, and thus controlling the flow rate to a prescribed level.\nHowever, in the case of measuring the flow rate of liquid CVD material by means of the liquid mass flow controller, a large error exists between the true and measured values because of much difference in specific heat and volume per unit time each of the CVD material flowing inside the piping as compared with the measurement of the flow rate of a gaseous CVD material with a gas mass flow controller. Hence it has been impossible to produce a semiconductor thin film with a very high accuracy of film thickness. On the other hand, a method in which the flow rate is measured with flow variable micropump involves a large error on account of pulsation flow of the liquid CVD material.\nIn such circumstances, an object of the present invention is to provide a vaporizing and supplying method for subjecting a liquid CVD material to flow rate control with a liquid flow rate controller, supplying a vaporizer with the material, vaporizing the same, and supplying a semiconductor manufacturing apparatus with the vaporized material, characterized in that the method enables a semiconductor thin film to be manufactured with an extremely high accuracy, while controlling the flow rate of the liquid CVD material with a high accuracy.\nOther objects of the present invention will become obvious from the text of this specification hereinafter disclosed.\nUnder such circumstances, intensive extensive research and development were accumulated by the present inventors in order to solve the above-described problems involved in the prior arts. As a result, the following has been found. Specifically, by installing in parallel, a plurality of liquid flow rate controllers, preferably at least two types of liquid flow rate controllers each having a different controllable range of flow rate and by using in combination so as to enable any of the aforesaid controllers to be used alone, or in plurality at the same time, or switchingly one after another, it is made possible to supply a vaporizer with the liquid CVD material in a wide range of capacity including a small to large capacity with a high accuracy, vaporize the material and supply a semiconductor manufacturing apparatus with the vaporized material. Thus, the present invention has been accomplished by the foregoing findings and information.\nThat is to say, the present invention is related to a vaporizing and supplying method for subjecting a liquid CVD material to flow rate control with a liquid flow rate controller, supplying a vaporizer with the material, vaporizing the same, and supplying a semiconductor manufacturing apparatus with the vaporized material, which comprises installing in parallel, a plurality of liquid flow rate controllers for the same type of liquid CVD material, and supplying the vaporizer with the material at a variable flow rate thereof by altering the single use of any of the controllers to simultaneous use of a plurality thereof or vice versa, and/or switching any of the controllers one after another."}
{"text": "Currently, most of price tags, price shelf labels, tickets, mass transportation passes and small signage are printed on paper. The disadvantage of displaying information on paper obviously is the consumption of paper, which not only is an environmental issue, but also a cost concern as the paper with information printed on must be discarded or recycled when the information needs to be updated. In addition, for certain time-sensitive systems, such as a pricing label or tag system, using paper to display information is labor intensive and often causes time delays and errors.\nThe concept of an electronic display device may be applied to these products. However, the high cost associated with an electronic display device prevents it from being widely used, especially for signs, labels, tickets and price tags which have to be regularly updated."}
{"text": "1. Field of the Invention\nThis invention relates generally to construction tools, and more particularly to a tool used to brace a masonry or other walls during construction.\n2. Description of the Related Art\nIt is well known that until a masonry wall is permanently supported, it is unstable and can fall, especially if wind loads on it are substantial. Thus, during the construction process braces and other devices are used to support the walls until the building is completed or the walls are at least self-supporting. Many devices used for supporting masonry walls are cumbersome and heavy, and they only adjust to walls that are certain heights, require forklifts or several laborers to install, require extra space to store and use parts that are only useful for one particular purpose. Thus, when a wall supporting apparatus is being installed, there is very little flexibility as to how these braces can be installed.\nAdditionally, existing wall bracing devices are made in certain sizes that are not flexible for use by contractors. Contractors must purchase different length braces for different wall heights. The bracing devices also have end pieces that are fixed to the poles, thereby making the bracing device suitable only to install after the wall is built and not while the wall is being constructed. Furthermore, since most bracing devices are long enough to be installed at high points in the wall, such as heights above 10′, they are conventionally too long and cumbersome to be transported easily.\nThe need exists for a wall bracing apparatus with parts that are easily adaptable to any height and for use in many circumstances, thereby eliminating searching effort, simplifying the installation process, and making the apparatus easily transported and stored when the apparatus is not being used."}
{"text": "The invention relates to a wheel suspension system for steerable rear wheels of a motor vehicle.\nRear wheel suspension systems of the generic type are known. For example, German Patent Document DE-OS 36 30 984, discloses an arrangement in which an arm connected at one end with the wheel carrier extends in the longitudinal direction of the vehicle, and is connected at its other end with a hydraulically operable adjusting element. The wheel carrier is pivotally supported between joints of transversely aligned control arms which are held on the vehicle body side. In such an arrangement, because of the arm connected with the wheel carrier, the position of the adjusting element is fixed with respect to the vehicle, and can be changed only at a relatively high cost.\nDue to space limitations (and also for kinematic and elastokinematic reasons), it may be necessary, particularly in the case of a double wishbone axle having an additional separate rod link, that the adjusting element not be connected directly with the wheel carrier.\nIt is therefore an object of the present invention to provide a wheel suspension for steerable rear wheels of motor vehicles in which direct control of a wheel suspension link is made possible for achieving a rear-wheel locking angle. This object is achieved according to the invention, in a double-wishbone suspension system wherein the wheels are adapted to react resiliently as a result of longitudinal or lateral forces and which, at the same time, can be controlled for the adjustment of a rear wheel steering angle by way of an adjusting element, such as a hydraulic control cylinder, a motor operator or similar devices. The two wishbones arranged in planes situated above one another, are pivotally connected to the wheel carrier as well as to the vehicle body and, either in the lower plane or in the upper plane, are supplemented by an additional rod link. The arrangement of the wheel suspension links and bearing members is such that a resilient rotating axis is established outside the wheel center longitudinal plane and-with respect to the driving direction--behind the wheel center transverse plane. (The words \"front\", \"rear\" and \"behind\" are used herein relative to the forward driving direction F as shown in FIG. 1.) Under the influence of outside forces and torques in the wheel contact area, the wheel rotates around this axis in the wheel contact surface. The extent of the rotation, for example in the case of a bump toe-in, is essentially a function of the elasticity of a control bearing between the adjusting element and a control arm, and the elasticity of the additional rod link. The control bearing also has an advantageous effect on the active adjustment of the rear-wheel locking angle by way of the adjusting element and, for this purpose, has a correspondingly defined marking.\nDepending on the construction of the double-wishbone axle and the connection of the adjusting element with one of the control arms, a corresponding active center of rotation for the wheel is disposed either in the intersecting point of the longitudinal axes of the control arms of the wheel suspension links or directly in the wheel-carrier-side joint of the wheel suspension link. As noted previously, the passive center of rotation of the wheel suspension system is determined by the elastic swivel axis and by the wishbones.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "This invention relates to wireless data communications networks, and in particular to arrangements for communications between mobile data handling units and a central computer using wireless data communications.\nThe assignee of the present invention supplies a wireless data communications system known as the Spectrum 24 System, which follows the radio data communications protocol of IEEE Standard 802.11. In the system as implemented, mobile units are in data communication with a central computer through access points. The access points may communicate with a central computer or computers over a wired network. Each of the mobile units associates itself with one of the access points. The access points in this system are functional to perform all the implemented requirements of the standard protocol, including, association and roaming functions, packet formulation and parsing, packet fragmentation and re-assembly encryption and system access control. In order to maintain order and reduce radio communications each access point must determine which of the data communications received over the wired network from the central computer is destined for a mobile unit associated with that particular access point. This requirement adds significant computational capacity to the access point, increasing the cost thereof.\nIn addition, in applications that must support a high volume of data communications from multiple users, such as systems supporting a self-service shopping system, hospital systems, systems that include paging or voice data links to many users, or systems supporting communicating with electronic shelf labels, additional access points are required to support the data communications traffic, increasing the overall system cost.\nThe cost of an operational access point is dependent not only on the complexity thereof and the requirement for high speed processing of data pockets for purposes of selecting those destined for mobile units associated with an access point, but the additional cost of the installation of electrical power to the location of the access point, and the cost of a power supply to convert AC electrical power to DC power for the circuits of the access point. Further cost may be involved in physically mounting the access point hardware and antenna.\nIn prior systems each access point is connected on an Ethernet wired network to the central computer. The access points are required to determine the identity of mobile units which have become associated with them and to extract from the data packets on the Ethernet network those packets addressed to a mobile unit associated with the access point. This requirement has led to significant processing burden for the access points and led to increased cost for the access points.\nIn the system described in my prior published International Patent Application WO 099 37047, published Jul. 22, 1999, the central computer communicates over an Ethernet wired network with an intelligent switching hub. Alternately a token ring network can be used. The switching hub determines the destination of each packet and routes packets to an access point if the destination of the packet is a mobile unit associated with the access point. To achieve this function, the hub is an intelligent hub which maintains a routing list of mobile units and their associated access point according to the port of the hub.\nIn practice, the hub need only maintain a source list for those access points connected to the hub and mobile units associated with the access points connected to the hub. Thus, if a packet is received at a hub over the Ethernet with a destination address which is not associated with that hub, the packet is ignored. The hub will route the packet to an access point only if the destination address of the packet is identified on the list. When a packet is received on a hub port associated with a communications line connected to an access point, the source address is associated with the hub port in the list. The packet is routed either to the Ethernet connection or to another port according to the destination address.\nBy determining destination address in the hub and maintaining the association of a mobile unit address with an access point connected to a port of the hub in a routing list of the hub, the functionality required of the access points is greatly reduced. The access point acts merely as a conduit sending RF transmissions of packets received on its communication line, and receiving transmissions from associated mobile units and providing Ethernet packets to the hub. In addition, the access point must provide mobile unit association functions and other 802.11 protocol functions, as provided in the Spectrum 24 system, and may also provide proxy polling responses for associated mobile units that are in power saving mode.\nThe prior system may have a large number of access points, each with a memory containing program instructions for carrying out the various required functions. This distribution of processing makes it difficult to upgrade a system or to provide changes in system configuration because any upgrade or change may require changes to the program code in each of the access points. Such distribution of processing functions also makes system management functions, such as load balancing or access control more difficult.\nIt is therefore an object of the present invention to provide an improved wireless data communications methods and systems having lower cost, to enable the economical provision of reliable wireless data communications with increased capacity in complex installations or at reasonable cost or simple installations."}
{"text": "The disclosed technology relates to halogenated polymer compounds, such as chlorinated polyvinyl chloride (CPVC) compounds, being free of heavy metals, particularly heavy metals in the form of stabilizer additives, the most common heavy metals being lead, tin, zinc, cadmium, and barium.\nHalogen containing polymers tend to degrade or deteriorate when processed. Generally, the difference between the processing temperature and the degradation temperature is very small. Therefore, there is a risk that during the processing these halogen containing polymers will degrade. When such polymers degrade, it is believed that the halide acid generated by the polymer attacks the components of the processing equipment. Also, this acid further catalyzes de-hydrohalogenation reactions, such as HCl elimination, and additional degradation of the polymer.\nStabilizers have been developed to help deter such degradation. For example, heavy metal compounds such as tin are commonly used. However, heavy metal compounds are considered environmentally unfriendly, and there is a move to limit their use in many jurisdictions. In anticipation of increased regulation of heavy metal stabilizers in halogenated polymers, there is a desire in the industry to discover replacement stabilizers.\nZeolites and metal carboxylates have been taught as co-stabilizers in halogen containing polymers, generally in addition to heavy metal stabilizers, but also in combination with other types of stabilizers.\nZeolites are effective acid scavengers for halogen containing polymers and can enhance the polymer's thermal stability. Acid scavengers are compounds that react with acid groups in the polymer reaction mixture to form compounds that are typically chemically inert. Such acid scavengers are often required as co-stabilizers in halogenated polymers to neutralize the acid generated during melt processing. However, the use of zeolites as acid scavengers in halogen containing polymer compounds has been limited for several reasons. First, the zeolites generally have a large particle size, generally in the range of about 3 to about 6 microns. The large size of the zeolite particles is thought to not only cause surface blemishes on the finish of the end product made from such a polymer but also to diminish the physical properties, for example impact strength, of such polymers. Further, outgassing occurs frequently with polymers containing zeolites due to the evolution of water from the zeolite from the heat generated during melt processing. Such outgassing can cause foaming and affect the physical properties of the final product.\nMetal carboxylates are also mentioned as potential co-stabilizers. However, when included at too high a concentration metal carboxylates have a tendency to plate-out (separate from the halogen polymer containing compound).\nU.S. Pat. No. 4,371,656 describes an “improved” metal substituted zeolite covered by an organic substance and having a water content of 8 weight percent (wt %) or less for use as a stabilizer for halogen containing resins. The zeolite is taught to be suitable as it is, but that it may function more effectively when used in combination with other stabilizers.\nA number of references teach the use of zeolite as a stabilizer in combination with other stabilizers. For example, U.S. Pat. No. 5,582,873 discloses a composition comprising a halogen containing polymer, a zeolite as the acid scavenger and a heat stabilizer selected from the group consisting of mixed metal stabilizers (such as calcium/zinc or calcium/barium), organotin stabilizers, lead stabilizers, metal free stabilizers or any combination thereof. Similarly, U.S. Pat. No. 6,414,071, to Wypart et al., issued Jul. 2, 2002 teaches a stabilized halogenated polymer wherein the stabilizer is a synthetic crystalline aluminosilicate of formula M2/nO.Al2O3.ySiO2.wH2O. For CPVC resins, the patent teaches and exemplifies the synthetic crystalline aluminosilicate in use with heavy metal stabilizers, including tin.\nOther references teach the use of zeolite in combination with sodium carboxylate. However, no references have been found that do not include other stabilizers with the zeolite/sodium carboxylate combination. For example, EP2083044A1 to Reagens S.p.A. teaches a composition for stabilizing halogen-containing polymers comprising disodium adipate and a compound selected from the group consisting of M(ClO4) and (CF3SO3)nM. Zeolite and carboxylates are also taught as additives that can optionally be included in the EP′044 composition. However, the patent discourages the use of zeolites by noting that zeolites tend to release water, which prevents their use as stabilizers in some rigid or semi-rigid PVC applications. The patent does not teach or exemplify a formulation containing CPVC, a zeolite and a carboxylate. Similarly, EP2363431A1 to Reagens S.p.A. teaches a composition for stabilizing halogen-containing polymers comprising an alkali metal salt of adipic acid and an alkali metal salt of a carboxylic acid. Zeolite is taught as an optional additional additive in the EP′431 composition, in an amount from 0.1 to 3 parts by weight per 100 parts by weight of polymer. The patent does not teach or exemplify a formulation containing CPVC and a stabilizer package consisting of a zeolite and a carboxylate.\nIt would be beneficial to the industry to prepare an inexpensive and readily available alternative to current stabilizer systems for halogenated polymers, such as CPVC resins, without the need for a heavy metal stabilizer."}
{"text": "The popularity of commercial wireless communications services (e.g., wireless telephony, wireless network access, and wireless email) has substantially increased during recent years. In many cases, users, such as consumers, mobile workers, emergency response personnel, and/or the like, now utilize these services for both personal and business communications. Likewise, users are increasingly relying on these services. For example, some households forgo wired telephone service in favor of wireless telephone service, some government agencies rely on these services for both routine and emergency communications, and businesses rely on these services to communicate with customers and mobile workers. Correspondingly, the cost (both financial and nonfinancial) of outages is also increasing.\nTypical commercial wireless communications service (CMRS) providers rely on remote facilities to facilitate the provision of services. For instance, CMRS providers rely on cell sites (e.g., base stations, radio repeaters, wireless to backhaul interfaces, etc.) to facilitate some communications services. If a cell site experiences a loss of commercially-provided electrical power, users near the cell site may experience a service outage. Power outages are an example of a common cause for cell site failures. For example, natural disasters, rolling brownouts, accidents, and/or the like may result in power outages. While most cell sites include some form of back-up power (e.g., generators and/or batteries), these forms of back-up power may not provide sufficient power during lengthy power outages and may require servicing, monitoring, and on-site maintenance. During lengthy power outages, use of commercial wireless communications services may increase due to users' needs and/or desires. Further, pending regulations may require commercial wireless communications service providers to provide cell sites with at least seven days of back-up power."}
{"text": "1. Field of the Invention\nThe invention relates generally to test systems, and more specifically, to test systems in the automatic test equipment (ATE) area.\n2. Discussion of the Related Art\nToday, most test equipment in the ATE area involves one or more computer systems. A computer system that implements testing steps must be reliable and must generate test steps in a deterministic way. That is, the computer systems must be capable of producing consistent results. One problem with such a computer is cost—design of a complete computing system suitable for use in testing (e.g., one that produces consistent results) requires not only the expensive hardware design and maintenance but also expensive software design involving operating system, user interface, drivers, tools, etc. More particularly, deterministic execution can be achieved using expensive computing hardware and software and/or distributed hardware within the tester, such as individual control processors within the instrument or other testing resource. As a result, computer systems suitable for ATE are prohibitively expensive for most companies to purchase and maintain.\nAutomated test equipment (ATE) is computer-controlled equipment that tests electronic devices for functionality and performance. ATE devices also conduct stress testing with minimal human interaction. A typical ATE tester includes control hardware, sensors, and software that collects and analyzes the test results. Such an ATE tester relies on particular test instruments to perform the testing and return results to the control software. Conventional ATE testers and methodologies are widely considered to only be cost efficient for high-volume testing.\nThis widely held belief stems, in part, from the design and operation of conventional ATE testers. As it stands now, a manufacturer can create a new device or system. A key element in the design and construction of any new device or system is the ability to test the device or system for quality, functionality, and serviceability. The first step in testing that new device is to determine if anyone already has a tester designed and working that may cover the new device. If one does not exist, the only option is to have a suitable tester built. The current marketplace takes a long time to determine if a new tester should be built, let alone to begin actual construction of such a tester. Determining whether or not a new tester covering the new device would have wide demand, or any at all beyond the single use, ultimately determines whether the new tester will be made. If it does not, that often ends the possibility of getting a new tester. If it does, the new tester still must be designed and fabricated. An average time for completing such a process is on the order of years. For instance, it is not uncommon for it to take two or more years to produce a new tester. In that time frame, most manufacturers are already looking to test a newer device, so such tester may be obsolete even before it is released."}
{"text": "Fibers with diameters less than a micron can be formed using electrospinning processes where a droplet of polymer solution or melt is elongated by a strong electrical field. The resulting fibers are collected as non-woven mats with extremely large surface to volume ratio; which are useful for various applications including filtration. Most previous studies on electrospinning have focused on fibers from polymer solutions, i.e., are directed to solution electrospinning. Current solution electrospinning apparatus and processes have the disadvantages of requiring a dissolving step, of requiring solvent recovery and disposal or complete recycling if the process is to be environmentally friendly, of having low production rates because of the dissolving and solvent recovery/recycling steps detracting from obtaining high throughput, of not being adaptable to polymers such as polyethylene, polypropylene, polyethylene terephthalate and polybutylene terephthalate, which are not dissolvable in acceptable solvents at room temperature, of requiring regulation of a plurality of parameters to adjust molecular properties and solidification and of requiring apparatus not readily provided by adaptation of conventional existing facilities for fiber/non-woven production for most polymers since these are based on melt treatment. Melt electrospinning apparatus and process which would avoid these disadvantages and provide useful production of fibers/non-wovens have not heretofore been developed. Moreover, no attempts have been made to provide solution electrospinning apparatus and processes which are suitable for operation on polymers which are not dissolvable in acceptable solvents at room temperature."}
{"text": "Pelargonium hortorum. \nxe2x80x98Clips Red IIxe2x80x99.\nThe present invention comprises a new and distinct cultivar of geranium, botanically known as Pelargonium hortorum Bailey, and hereinafter referred to by the cultivar name xe2x80x98Clips Red IIxe2x80x99. The new cultivar is propagated from a seedling resulting from the cross of an unnamed and unpatented proprietary line known as 7108-1 which is a red single flowered geranium and an unnamed and unpatented proprietary line known as 6087-4 which is a red semi-double geranium flower.\nxe2x80x98Clips Red II xe2x80x99is a product of a planned breeding program intended to create new geranium cultivars with red colored flowers, semi-double flower form, dark green foliage, vigorous growth and superior cutting productivity.\nThe new cultivar was created in 1997 in Gilroy, Calif. and has been repeatedly asexually reproduced by cuttings in Gilroy, Calif. and Guatemala over a three year period. It has also been trialed at Gilroy, Calif. and Andijk, The Netherlands. It has been found to retain its distinctive characteristics through successive propagations; and this novelty is firmly fixed."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming system comprised of an image forming apparatus that forms an image on a sheet, and a clear coating apparatus connected to the image forming apparatus on a downstream side of the same so as to perform clear coating on the sheet using clear toner which becomes clear (transparent) after being fixed on the sheet.\n2. Description of the Related Art\nIn recent years, there has been an increasing demand for full-color electrophotographic image forming apparatuses ensuring high image quality. Further, it has been increasingly desired to realize a high-quality photographic tone image mode in an electrophotographic image forming apparatus, with proliferation of digital cameras as a major reason.\nFurther, there has been proposed an image forming system that is comprised of an image forming apparatus and a clear coating apparatus continuously connected thereto, and is configured to form a high-quality photographic tone image on a sheet, and then place thereon a toner image which becomes transparent when fixed, to thereby perform clear coating on the formed image (see e.g. Japanese Laid-Open Patent Publication (Kokai) No. H03-13079).\nFurther, there has been a system including apparatuses continuously connected to each other, in which an upstream one of the continuously connected apparatuses performs printing on a first surface of a sheet, and a downstream one of the apparatuses performs printing on a second surface of the sheet, to thereby realize double-sided printing without the provision of a sheet-inverting mechanism (see e.g. Japanese Laid-Open Patent Publication (Kokai) No. 2003-207954).\nThe above system with the continuously connected apparatuses is aimed at improving productivity, and hence double-sided printing is not performed within any of the apparatuses. Therefore, when the upstream one of the apparatuses performs double-sided printing on each sheet and discharges the sheet, according to a conventional job sequence, if the downstream one of the apparatuses also perform double-sided printing on the sheet, this causes the problem of an improper page order of the discharged printed sheets."}
{"text": "Not applicable.\nNot applicable.\n1. Field of the Invention\nThis invention relates to the genes and enzymes involved in cell wall synthesis in bacteria, and particularly to the inhibition of such enzymes.\n2. Background of the Invention\nThe molecular target of many naturally-occurring antibiotics, including fosfomycin, cycloserine and b-lactams, is the synthesis of the bacterial cell wall. The frequency with which these types of antibiotics arose in evolution indicates that the pathway of cell wall biosynthesis is a particularly effective point of attack against bacteria. Genetic studies confirm the soundness of this process as a target, as temperature-sensitive alleles of the intracellular pathway genes are lytic, and therefore lethal. Since the building blocks of the cell wall are highly conserved structures in both Gram-positive and Gram-negative bacteria, but are unique to the eubacteria, novel inhibitors of cell wall formation are expected to be both broad spectrum and safe antibiotics.\nThe bacterial cell wall is a polymerxe2x80x94a single molecule composed of peptidoglycanxe2x80x94that defines the boundary and shape of the cell. Assembled by crosslinking glycan chains with short peptide bridges (Rogers, H. J., H. R. Perkins, and J. B. Ward, 1980, Biosynthesis of peptidoglycan. p. 239-297. In Microbial cell walls and membranes. Chapman and Hall Ltd. London), the completed structure is strong enough to maintain cell integrity against an osmotic pressure differential of over four atmospheres, but also flexible enough to allow the cell to move, grow and divide.\nThe construction of the peptidoglycan begins in the cytoplasm with an activated sugar molecule, UDP-N-acetylglucosamine. After two reactions (catalyzed by MurA and MurB) that result in the placement of a lactyl group on the 3-OH of the glucosamine moiety, a series of ATP-dependent amino acid ligases (MurC, -D, -E, and -F) catalyze the stepwise synthesis of the pentapeptide sidechain using the newly synthesized lactyl carboxylate as the first acceptor site. After attachment of the sugar pentapeptide to a lipid carrier in the plasma membrane, another glucosamine unit is added to the 4-OH of the muramic acid moiety. The completed monomeric building block is moved across the membrane into the periplasm where the penicillin-binding proteins enzymatically add it into the growing cell wall.\nBecause the pentapeptide sidechain is not synthesized ribosomally it contains more diverse chemical functionality than a typical peptide, both structurally and stereochemically. Two of the enzymes catalyze the addition of D-amino acids (MurD and MurF) and MurE mediates the formation of a peptide bond between the g-carboxylate of D-glutamate and the amino group of L-lysine. Presumably these structures render the exposed peptidoglycan resistant to the action of proteases, but they also imply that the active sites of the enzymes must have unusual structures in order to handle the somewhat uncommon substrates. These unusual active sites are targets to bind novel inhibitors that can have antimicrobial activity.\nAmong these potential enzyme targets is MurD. The first partial purification and characterization of a D-glutamate-adding enzyme was from Staphlococcus aureus (Ito, E. and J. L. Strominger, 1962, Enzymatic synthesis of the peptide in bacterial uridine nucleotides: Enzymatic addition of L-alanine, D-glutamic acid, and L-lysine. J. Biol. Chem. 237: 2689-2695; Nathenson, S. G., J. L. Strominger, and E. Ito, 1964, Enzymatic synthesis of the peptide in bacterial uridine nucleotides: purification and properties of D-Glutamic acid-adding enzyme, J. Biol. Chem. 239: 1773-177), followed by studies in more detail on the isolated E. coli enzyme (Blanot, D., A. Kretsovali, M. Abo-Ghalia, D. Mengin-Lecreulx, and J. van Heijenoort, 1983. Synthesis of analogues of precusors of bacterial peptidoglycan. p. 311-314. In Peptides. Blaha, K. and P. Malon, eds. pp. 311-314, Walter de Gryter and Co. Berlin, N.Y.; Jin, H., Emanuele, J. J., Jr., Fairman, R., Robertson, J. G., Hail, M. E., Ho, H.-T., Falk, P. and Villafranca, J. J., 1996. Structural studies of Escherichia coli UDP-N-acetylmuramate: L-alanine ligase. Biochemistry 35: 14423-14431; Ito E. and J. L. Strominger, 1973. Enzymatic synthesis of the peptide in bacterial uridine nucleotides: Comparative biochemistry. J. Biol. Chem. 248: 3131-3136; Michaud, C. D. Blanot, B. Flouret, and J. van Heijenoort, 1987. Partial purification and specificity studies of the D-glutamate-adding and D-alanyl-D-alanine-adding enzymes from Escherichia coli K12. Eur. J. Biochem. 166: 631-637). Recently, a purified recombinant E. coli MurD was reported (Pratviel-Sosa F, D. Mengin-Lecreulx and J. van Heijenoort, 1991. Over-production, purification and properties of the uridine diphosphate N-acetylmuramoyl-L-alanine:D-glutamate ligase from Escherichia coli. Eur. J. Biochem. 202 (3):1169-1176) and genes encoding MurD have been cloned from several species of bacteria including E. coli (Ikeda, M., M. Wachi, F. Ishino, and M. Matsuhashi, 1990, Nucleotide sequence involving murD and an open reading frame ORF-Y spacing murF and ftsW in Escherichia coli. Nucleic Acids Res. 18:1058; Mengin-Lecreulx, D., C Parquet, L. Desviat, J. Pla, B. Flouret, J. Ayala and J. van Heijenoort, 1989, Organization of the murE-murG region of Escherichia coli: Identification of the murD gene encoding the D-glutamic-acid-adding enzyme. J. Bacteriol. 171: 6126-6134) and B. subtilis (Daniel, R. A., and J. Errington, 1993, DNA sequence of the murE-murD region of Bacillus subtilis 168. J. Gen. Microbiol. 139:361-370; Henriques, A. O. de Lencaster, H. and P. J. Piggot, 1992, A Bacillus subtilis morphogene cluster that includes spoVE is homologous to the mra region of Escherichia coli. Biochimie. 74: 735-748). Compounds have been designed and synthesized that have inhibitory activity against the E. coli enzyme (Tanner, M. E., S. Vaganay, van Heijenoort, J., and D. Blanot, 1996, Phosphinate Inhibitors of the D-Glutamic Acid-Adding Enzyme of Peptidoglycan Biosynthesis. J. Org. Chem. 61: 1756-1760), although they do not have antibacterial activity.\nPolynucleotides and polypeptides of Streptococcus pyogenes MurD, an enzyme involved in bacterial cell wall biosynthesis are provided. The recombinant MurD enzyme is catalytically active in ATP-dependent D-glutamate addition reactions. The enzyme is used in in vitro assays to screen for antibacterial compounds that target cell wall biosynthesis. The invention includes the purified polynucleotides, purified enzymes encoded by the polynucleotides, and host cells expressing the recombinant enzyme and their use in assays."}
{"text": "The present invention is generally directed toward vibratory finishing machines and more particularly, toward a new and improved finishing machine of the type adapted to be used with steel finishing media.\nOne of the objectionable characteristics of utilizing steel finishing media in vibratory finishing machines is the difficulty in obtaining the desired \"flow\" of the media. One of the significant reasons for such difficulty resides in the need for utilizing certain chemical compound additives in the finishing operation which have been found to severely hinder or preclude media flow, which in turn results in precluding effective workpiece finishing. Although satisfactory media flow has been obtainable in the past through the use of large counterweights, such an approach has not been feasible with steel media because of the high magnitude of media-workpiece impact which prevents obtaining the fine finish desired for the type of parts that are traditionally processed with steel finishing media.\nThe present invention provides a new and improved vibratory finishing machine which has been found to provide highly improved media flow at low magnitudes of media-workpiece impact, thereby obviating the various objectionable characteristics of the prior art machines. The finishing machine of the present invention utilizes a unique combination of features, including the provision of support springs for the finishing tub at a position adjacent the peripheral rim or edge of the tub, the use of an asymmetrically located weighted vibratory shaft and the additional use of an asymmetrically located tub weight that is preferably located on the opposite side of the longitudinal axis of the tub from the vibratory shaft, as will hereinafter be described in detail.\nIt is accordingly a general object of the present invention to provide a new and improved vibratory finishing machine.\nIt is a more particular object of the present invention to provide a new and improved vibratory finishing machine that is particularly, although not necessarily, adapted for use with steel finishing media and is intended to provide improved finishing action, as compared to prior known finishing machines.\nIt is another object of the present invention to provide a new and improved finishing machine of the above described type which incorporates a plurality of suspension or support springs that are located above the center of gravity of the finishing tub, and specifically are located directly beneath the peripheral flange or rim of the tub.\nIt is still another object of the present invention to provide a new and improved finishing machine that is provided with an asymmetrically located weighted vibratory shaft.\nIt is a further object of the present invention to provide a new and improved vibratory finishing machine that may be provided with an asymmetrically located tub weight which is preferably located on the laterally opposite side of the longitudinal tub axis from the weighted vibratory shaft.\nIt is still a further object of the present invention to provide a new and improved vibratory finishing machine wherein improved media flow is achieved in the range of low media impact with the media workpieces, whereby to improve the operational performance of the machine of the present invention.\nOther objects and advantages of the present invention will become apparent from the following detailed description taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates to an image recording apparatus with a desired exposure device mounted thereon as an option.\nAs an image recording apparatus of the type using a photosensitive pressure-sensitive recording medium carrying thereon an immense number of pressure-rupturable, various kinds of image recording apparatuses have been known in the art which include a full-color copying machine, a CRT printer, a film printer. These machines or devices have been designed to have a relevant handling quality and requisite functions depending on the intended use thereof.\nHowever, notwithstanding the fact that these devices have a basic arrangement common to all these devices, the final arrangements of these devices are considerably different due to a difference in intended use. To design and build up each of the different types of image recording apparatuses, a great deal of time and expense has been required."}
{"text": "There is a need in the field of microbial metabolic engineering for tunable promoters and novel regulatory switches for the inducible expression of heterologous proteins. The “old standard” lac promoter/Lac repressor system is still widely used. However, the commonly used inducer, isopropylthio-beta-D-galactoside (“IPTG”), is expensive and thus not practical for large-scale bioprocesses. Additionally, higher concentrations of IPTG increase the metabolic burden on the cell, in turn reducing the maximal expression of the target gene (Donovan et al., J. Ind. Microbiol. 16:145-154 (1996)). The few available alternatives also have limitations.\nOne possible solution to this problem is to appropriate existing genetic systems of transcriptional regulators to enhance heterologous gene expression. A number a of transcriptional regulators are known. For example, the LysR family of transcriptional regulators is one of the largest groups of transcriptional regulators in prokaryotes (Schell, Annu. Rev. Microbiol. 47:597-626 (1993)). Currently, there are over 80 known members of this regulator family. Proteins having greater than 20% amino acid identity with another LysR family member or having the consensus sequence of the N-terminal region of the LysR family are considered to be members of this regulator family. LysR family members are also commonly found in the size range of 276 to 324 amino acids, bind to similar DNA sequences in the absence of inducers, have promoters that are located close to or overlapping those of the regulated target gene, and most can repress their own transcriptional levels 3- to 10-fold. Activation of the regulated target gene occurs in the presence of inducer and usually results in a 6- to 200-fold increase in regulated target gene transcription. Regulated target genes are diverse and have numerous functions.\nRecently, the gene encoded by open reading frame (“ORF”) b3243 in Escherichia coli (“E. coli”) has been demonstrated to function via quorum sensing (Sperandio et al., Infect. Immun. 70:3085-3093 (2002)). Quorum sensing is the ability of bacteria cells to communicate with one another through perception of the accumulation of signaling molecules based on bacteria cell density. The gene encoded by ORF b3243 is up-regulated via quorum sensing resulting in a 23-fold increase in transcription of the gene. The protein produced by the b3243 ORF was not purified. The gene was found to have a role in the regulation of the LEE genes involved in a type III secretion system, a pathogenecity system that serves to translocate, upon contact with eukaryotic host cells, proteins from the bacteria cytoplasm into the host cell cytoplasm. As the b3243 ORF gene is a putative regulator of the LysR family, the b3243 ORF gene itself was able to induce a four-fold induction of LEE1 transcription. No inducer of the b3243 ORF gene was identified. One in the art will appreciate that this inefficient induction prevents the b3243 ORF gene from being a viable promoter/regulator system without the identification of its inducer.\nThe use of promoter/reporter gene constructs is well known in the art (Serebriiskii and Golemis, Anal. Biochem. 285:1-15 (2000); Spergel et al., Prog. Neurobiol. 63:673-686 (2001); Yarranton, Curr. Opin. Biotechnol. 3:506-511 (1992)). Particularly, reporter systems utilizing β-galactosidase, green florescent protein, luciferase, and chloramphenicol acetyl transferase (CAT) are all commonly used in the art. Additionally, reporter systems such as luxCDABE are useful because this operon contains all of the genes required for bioluminescent reporting ((Van Dyk et al., Proc. Nat. Acad. Sci. USA 98:2555-2560 (2001)).\nThe luxCDABE operon has been utilized to create a collection of random gene fusions, comprising 27% of the known or predicted transcriptional units of E. coli (Van Dyk et al., J. Bacteriol. 183:5496-5505 (2001)). Treatment of E. coli cells containing these gene fusions with nalidixic acid, a quinolone, results in selective up-regulation of ten genes. Some of these up-regulated genes are LexA-regulated SOS genes, while others are not generally induced by DNA damage.\nAromatic compounds such as aromatic carboxylic acids are usually toxic to microorganisms. Numerous bacterial strains resistant to aromatic compounds, however, are known in the prior art (Diaz et al., Microbiol. Mol. Biol. Rev. 65:523-569 (2001)). Additionally, a LysR family member from Acinetobacter, BenM, is responsive to synergistic induction by benzoic acid, an aromatic carboxylic acid, and muconic acid (Bundy et al., Proc. Natl. Acad. Sci. USA 99:7693-7698 (2002)). Benzoic acid alone, however, produces minimal, if any, induction of BenM activity. Even the synergistic response with muconic acid produces only a four-fold increase in BenM activity.\nU.S. Pat. No. 5,292,643 issued to Shibano et al. on Mar. 8, 1994 describes genes related to fusaric acid resistance in variety of microorganisms. Specifically, genes capable of decomposing or detoxifying fusaric acid are disclosed. One of the genes postulated to be involved in fusaric acid resistance, fusB, shares some homology with the putativTe efflux transporter (PET) yhcP gene (Paulsen et al., FEMS Microbiol. Lett. 156:1-8 (1997)). Applicants incorporate by reference the co-owned and concurrently filed application entitled “PET Family of Efflux Proteins”, U.S. Patent Application No. 60/440,760, which describes new proteins efflux proteins whose expression may alter the expression of carboxylic acids.\nThe problem to be solved therefore is to discover facile and inexpensive methods of inducible expression of heterologous genes. Applicants have solved the stated problem through the discovery that the promoter elements of the yhcRQP operon are responsive to the expression of the yhcS regulator (a member of the LysR family of transcriptional regulators) whose expression may be induced by an inexpensive cadre of aromatic carboxylic acid inducers."}
{"text": "Affective computing involves the study and development of systems and devices that can recognize and analyze the expression of human emotions (affects). Such systems often rely on statistical models created from data collected by monitoring users and their affective response to various stimuli. A user's affective response to stimuli is typically a temporary departure from a steady baseline affective state. The baseline state is the state the user would typically be in when not exposed to any particular stimuli. In order to accurately model and understand the effects of various stimuli on the user, the user's affective response needs to be evaluated relative to the user's affective baseline value; otherwise, it is difficult to determine for instance, if a change in the user's affective state is due the stimuli, or is simply a natural return to the user's baseline affective state.\nOne common method to compute a baseline value is by measuring the user in a rested state before the exposure to stimuli (for example before an experiment starts). Another common method to compute a baseline is by weighting and processing the user's affective response values collected in a relatively wide window, for example by taking the average of values collected over a one hour. The problem with the former method for computing a user's baseline is that it is usually only applicable in controlled situations, for example, experiments in laboratories. In real world scenarios, it is likely the user will have neither the willingness nor the opportunity to provide an accurate baseline measurement from a calm period, before each exposure to stimuli. The later method described for computing baselines has a weakness, which can often make the baseline values computed by this method very inaccurate. This method does not take into account the situation the user is in. In real world scenario, at different times, the user may be in very different situations and the user's baseline value can change dramatically depending on the situation. For example, a user's affective baseline while working in the office might be quite different from the user's baseline when relaxing at home. In addition, the user's affective baseline can change dramatically in a very short when the situation changes. For example, a user's affective baseline state can change within seconds from something that can be described as calm and relaxed, for instance while cruising down highway 101, to tense and agitated when the user notices a police car in the rear-view mirror (along with the realization that the user was driving at a speed of 85 mph). Since the baseline value computed with a wide window changes slowly, it cannot adjust to a user's situation dependent affective baseline. This may make any models or predictions of the user's response to stimuli inaccurate."}
{"text": "Control of the temperature of reactions within an acceptable range has been widely investigated and the chemical industry has devised several arrangements, those commonly used being discussed in standard references and texts, e.g. one might consider the general teachings by Octave LEVENSPIEL in Chapter 19 of Chemical Reaction Engineering (published by John Wiley & Sons).\nThe prior art includes a conventional reactor designed to offer more control over the reactant temperature and this is known as the staged adiabatic packed bed reactor. This system uses an arrangement wherein a number of discrete, spaced apart zones of reaction are provided with means therebetween to control the temperature of the products leaving a first zone of reaction prior to entering the next reaction zone. No heat exchanging means is provided to control the temperature of the reaction in the zones of the reaction. Thus the reactant fluid entering the reactor at a desired temperature passes through a packed bed containing catalyst. Upon exiting this first stage, the reactant gas and any products will have a temperature higher or lower than that of the initial temperature depending upon the reaction thermal characteristics. A heat exchanger then heats or cools the reactant gas to a second desired temperature, which may or may not be equivalent to the temperature of the first, before passing to the next packed bed i.e. the second stage. This sequence is repeated until the desired conversion is obtained. Thus the temperature profile of the reaction will be stepped within an acceptable range of temperature, and will therefore not be truly isothermal.\nThe preferred heat exchanger panel for the purposes of the invention is one formed from a plurality of plates superposed and diffusion bonded to form a stack of plates, wherein fluid passages are defined in said stack by virtue of a pre-treatment of said plates wherein each plate is selectively configured to provide channeled or blank surfaces according to the desired pattern of fluid passages by a treatment to remove surface material e.g. by chemical etching, or hydraulic milling, or the like process to a desired depth. Optionally the chemical treatment may be augmented by a mechanical treatment using a suitable tool.\nSuch a pre-treatment of the plates is conducted in a manner analogous to manufacture of printed circuit boards (PCBs) and for this reason the heat exchanger design can be described as a printed circuit heat exchanger (PCHE). The application of the diffusion bonding technique for metal plates is well understood in the art of metal working and is applied for a variety of purposes e.g. in medical prosthesis manufacture.\nThis design of the PCHE has been proven by the designers of the proposed PCR system since 1985 when these compact heat exchangers were first introduced.\nA PCR type of reactor was designed by the present applicants and is the subject of a separate patent application (Ref:32 46271 WO-). Such a reactor is formed to provide at least one reaction zone, bounded by a heat exchanger formed from a plurality of plates superposed and diffusion bonded to form a stack of plates, wherein fluid channels are defined in said stack by virtue of pre-treatment of said plates wherein each plate is selectively configured according to the desired pattern of channels by a chemical treatment to remove surface material e.g. by chemical etching, to a desired depth. The fluid channels defined in the stack provide the opportunity to arrange for various reactant fluids to be conveyed in channels arranged in a heat exchange relationship with discrete channels containing at least one auxiliary fluid for controlling the temperature of the reactants.\nConsidering the example of a known ammonia converter, for a given ammonia content in the reactants there is a temperature for which the desired reaction rate is at a maximum. This is because the rate of synthesis of ammonia is the net result of the competing rates of the forward and reverse reactions. Consequently, by monitoring and controlling the temperature, it is possible to determine a temperature that favours the forward production reaction more than the reverse product dissociation reaction. In fact it is found that whilst increased temperature generally causes an increase in reaction, and indeed the desired forward product formation reaction rate increases favourably with temperature, at a certain approach to equilibrium conditions the concurrent increased rate of the reverse reaction begins to dominate, and to slow the overall synthesis rate. Maximum conversion in a bed of given size is therefore achieved if conditions remain on the maximum rate line. This is shown graphically in FIG. 1.\nAmmonia synthesis typically takes place at high pressures, greater than 100 bar and therefore creating bed volume is relatively expensive. Furthermore, the catalyst itself is costly. Consequently, an important element of ammonia synthesis reactor design is to attempt to maintain the bed conditions as closely as possible to the maximum reaction rate line, and hence to maximise the rate of ammonia synthesis in a bed of a given size.\nVarious approaches have been taken to achieving this result, mainly falling into the categories of:\n1. Quench-cooled multi-bed converters, in which cold reactant feed is injected between beds of catalyst,\n2. Tube-cooled converters, in which tubes carrying cold reactant are embedded in the catalyst bed, and\n3. Indirectly cooled multi-bed exchangers, in which heat is extracted from the hot reactants passing between beds in heat exchangers cooled by cold reactant feed.\nAn example of a quench type converter is disclosed in U.S. Pat. No. 3,663,179 wherein there is provided a vertically oriented container or reactor vessel, which is provided with an internal catalyst basket in which a bed or charge of catalyst particles is disposed. The basket is spaced from the container wall, and the feed fluid stream such as synthesis gas is passed onto the lower portion of the container and external to the basket. The feed fluid rises through the annular space between the basket and the container wall, and thus serves to cool the container and act as insulation against the thermal effects or hot spots generated within the catalyst bed. The rising warmed feed fluid stream is then heated to a suitable catalysis temperature by an internal heat exchanger. The hot fluid then flows to the catalyst bed in which the reaction takes place.\nA perforated pipe is disposed in the catalyst bed, and a cold quench fluid, which may be of a composition comparable to the feed stream is passed through the pipes and distributed into the hot reacting gas within the catalyst bed, to provide a cooling effect and thereby moderate the catalytic reaction.\nAmmonia synthesis is exothermic, typically operating in the 350° C. to 500° C. range, and the conversion factor is relatively low, typically less than 20%. As a result, it is possible for the cold feed stream to the reactor to be used to extract the reaction heat, at the same time preheating the cold feed stream to the required reaction temperatures. It is conventional practice to use the feed stream as a cooling medium in this way in ammonia converters. Related documents U.S. Pat. No. 4,230,669 and U.S. Pat. No. 4,230,680 describe such a converter in which the cold feed stream is used to extract the reaction heat and also includes a cold feed bypass line such that a controllable volume of the cold feed stream can by-pass the heat exchangers and therefore better control of the reactant temperature can be achieved.\nThis converter is an example of an indirectly cooled multi-bed exchanger as briefly outlined at point 3 above. It is an approach to bed cooling that is generally preferred by those in the art. In prior art converters though, only a very crude tracking of the maximum reaction rate curve is achievable in the reactor, due to the difficulties in cost-effectively arranging for more than two beds with inter-cooling by the feed stream. The general form of the temperature profile achieved with the common two-stage arrangement is illustrated in FIG. 2. The proposal outlined in U.S. Pat. Nos. 4,230,669 and 4,230,680 mentioned above actually has three catalyst beds, but even so performance leaves room for improvement and the design cannot be considered compact.\nOther examples of ammonia converters are described in the publication “Ammonia and Synthesis Gas, Recent and Energy-saving Processes” edited by F J Brykowski and published by Noyes Data Corporation in 1981.\nAn object of the present invention is to provide an improved reactor design, particularly one that is useful for the purposes of ammonia conversion and the like reactions, and is furthermore of a relatively compact design."}
{"text": "Dental restoration techniques include the use of polymers and other restorative materials to regain a tooth's function, form, and integrity. In a restorative technique referred to as photopolymerization, a polymer adhesive is applied to a tooth or portion thereof. The adhesive is then exposed to a light of a certain wavelength, which causes a chemical reaction to occur between the polymers of the adhesive. This reaction transforms the adhesive from a liquid or pliable state to a solid state intended to withstand strains, pressures, and other conditions often experienced by teeth."}
{"text": "The invention will be discussed generally with respect to deployment of a stent graft into the thoracic aorta but is not so limited and may apply to deployment into other body lumens.\nIt is important when delivering a stent graft by intraluminal or endovascular methods to know exactly where the device is in the vasculature particularly in relation to branch vessels which could be adversely occluded by placement of the stent graft. To avoid occlusion there have been proposed fenestrations along the length of a stent graft and scalloping at the proximal and distal ends of the stent graft to allow flow from a main vessel to a branch vessel through the fenestration or scallop.\nThe use of radiopaque markers around the periphery of the fenestration or scallop in this regard has been proposed but it would be advantageous if a more certain or reliable method could be devised to assist with placement.\nParticularly when deploying into a curved vessel it is known that a deployment device will take up a position in the vessel on the outside or greater radius side of the curved vessel and hence if a stent graft could be retained onto a deployment device so that a fenestration or scallop on the stent graft is in a selected position with respect to a portion of the deployment device when it is curved then more accurate positioning may be possible.\nAs a particular example, the aorta of a patient comprises an ascending aorta from the aortic heart valve, a thoracic arch and a descending aorta. Major branch vessels extend from the thoracic arch and occlusion of one or more of these upon placement of a stent graft in the thoracic arch could have serious consequences. The major vessels generally extend from the outside or greater radius side of the curved thoracic arch and so if a was stent graft retained onto a deployment device such that a scallop on the stent graft was on the outside of the deployment device when it is curved then correct placement would be surer.\nIt is known to retain the proximal and distal ends of a stent graft onto a deployment device to facilitate relative longitudinal and rotational movement of the ends but owing to tortuosity of vessels through which a deployment device must be progressed and the rotation necessity to achieve progression the final rotational position of a stent graft may not be fully known.\nIt is an object of this invention to provide a retention mechanism for a stent graft onto a deployment device for more accurate deployment of the stent graft therefrom or to at least provide a practitioner with a useful alternative.\nThroughout this specification the term distal with respect to a portion of the aorta, a deployment device or a stent graft is intended to mean the portion of the aorta, deployment device or stent graft further away in the direction of blood flow away from the heart and the term proximal is intended to mean the portion of the aorta, deployment device or end of the stent graft nearer to the heart. When applied to other vessels similar terms such as caudal and cranial should be understood."}
{"text": "The disclosure relates generally to deadlock avoidance in a multi-processor computer system with extended cache line locking.\nAs contemporary multi-processor computer designs have scaled to larger symmetric multiprocessing (SMP) systems, performance issues around contentious cache lines have become more noticeable to the system. In SMP systems, a traditional deadlock avoidance mechanism based on a length of time that an operation is valid for is implemented to address functional and performance issues. That is, the traditional deadlock avoidance mechanism is initiated to help an operation make forward progress, if that operation is valid for longer than a set time. Additionally, in SMP systems, a next instruction access intent (NIAI) lock instruction can be utilized by a processor core to purposefully hold onto a cache line for an extended period of time to allow the processor core to complete critical sections of code. In this case, while the NIAI lock instruction is purposefully holding the cache line, the traditional deadlock avoidance mechanism (which is based solely on the duration of an operation) can falsely activate and cause instruction sequences waiting on the hold to interfere with other unrelated instruction sequences."}
{"text": "This invention relates to corrosion protection systems, and more particularly to corrosion protection systems for joined dissimilar metals which are employed in an aqueous environment.\nPotable water, or similar type electrolyte normally contains certain corrosion causing components such as dissolved oxygen. Other oxidizing components of many varieties can be present in this medium, depending on the particular sample of water, the location from where the water was obtained as a geological function, certain water treatments such as softening, and the many other variable conditions to which water can be subjected.\nIt is common for potable water and other types of water systems to contain iron or iron bearing alloys which can form a galvanic corrosion cell between the iron bearing material and each oxidizing atom or molecule involved in a like reaction. These reactions can reach factors of 10.sup.(x) per square inch. Therefore, it has not been possible to render each separate corrosion cell inactive by direct action to each oxidizing atom or molecule.\nCathodic protection systems have been developed to reduce corrosion caused by oxidizing components found distributed in the potable or similar type water containing the vulnerable metal. In a cathodic protection system, the vulnerable metal is electrically charged by a variable or constant voltage in order to cause the formation of a monatomic hydrogen layer. This hydrogen layer, when undisturbed, forms an interface between the oxidizing components and the vulnerable metal. In a passive environment where little or no turbulence exists, as well as other factors which would have an effect upon the integrity of the monatomic hydrogen layer, this form of protection can be useful in the reduction of corrosion cell formation from the waterborne corrosion inducing components. Cathodic protection is even more effective when the base metal is coated, such as with glasslining, cement or epoxy, reducing the contact area between the vulnerable metal and the aqueous medium.\nHowever, impressed cathodic protection is not as effective and can actually accelerate corrosion in certain aggressive environments. Examples are an environment where there is turbulent water flow against the vulnerable metal removing the hydrogen layer and leaving a charged surface, or an environment characterized by elevated temperature, pH extremes, or contact of the vulnerable metal with a more noble dissimilar metal.\nFor dissimilar metal junctions, it can be shown that impressed cathodic protection of the vulnerable metal will accelerate the corrosion of the vulnerable metal by 20% to 45% or more when any significant turbulence is present. Naturally occurring cathodic protection that occures between joined dissimilar metals will cause only slight corrosion rate reductions in a low temperature environment, about 23.degree. Centigrade where very little to no turbulence is present, and the water sample is non-aggressive, tap water, untreated at point of use.\nThe major water heater manufacturers have long recognized the shortfall of cathodic protection against dissimilar metal junctions and consequently avoid the use of dissimilar metal junctions in commercial units. Many water heater tanks have di-electric inlet and outlet factory attachments. Thermostat copper and brass shanks, as well as relief valve probes, are plastic coated. Drain valves are plastic, instead of copper or brass whenever possible. Copper electric heating elements are tin coated. Also, the new class of heating elements are now commonly produced. These heating elements are comprised of stainless steel or an iron alloy called Incolloy which further reduces dissimilar metal corrosion. Plumbing codes almost always specify that di-electric unions must be used when copper pipe is run to and from a hot water heater.\nVarious methods have been proposed to prevent corrosion. For example, L. Applegate, Cathodic Protection, McGraw Hill, February 1960, pp 1-28, discloses a method to stop the reaction between separate bars of copper and iron which involves application of an impressed current to the dissimilar metals. However, this method as disclosed by Applegate only has a laboratory application. The metals are never directly joined. Furthermore, the positive current electrode is connected directly to the more noble copper bar which causes the copper to become very active and disintegrate. Moreover, Applegate discloses that the accepted conclusion is that dissimilar metals should not be used due to the unpredictable nature and widespread potential for system corrosion."}
{"text": "1. Field of the Invention\nThe present invention relates to a musical apparatus which automatically repeats a certain specified tune sequence of a musical tune or automatically starts playing from a specified point in a musical tune sequence.\n2. Technical Background\nAs is well known, there are many musical apparatuses which can automatically play a tune according to specified performance data. Among such apparatuses, there are some which can play all or part of the tune repeatedly, for example, Japanese Patent Application Laid-open Publication No. Sho 58-42093 discloses an apparatus which permits any part of a tune to be specified, and can repeatedly replay only the specified part of the tune.\nIn the apparatus shown in the above publication, the specified part is simply read out repeatedly, and there is no creativity in the manner of repetition, and the part of the musical tune is simply repeated many times over and over.\nOn the other hand, there is an automatic performance apparatus known to memorize a plurality of pattern data, and the order of playing the pattern data can be preprogrammed. However, this type of apparatus (pattern sequencer) could not repeatedly play any part of a tune, or start from any point in a tune."}
{"text": "1. Field of the Invention\nThe present invention relates to a display control apparatus and a camera system.\n2. Description of the Related Art\nIn recent years, as a monitoring camera system, a network camera which can freely perform a PTZR (pan, tilt, zoom, rotation) driving has been spread and many kinds of various functions have been provided. Therefore, there is a monitoring camera system in which video display units corresponding to a plurality of video sizes, functional dialogues (dialogue boxes) corresponding to a plurality of camera operations, and the like are displayed in a mixed manner.\nAs a system which handles a plurality of windows, a variety of systems such as general OS itself, image editing software, document editing software, remote communication tool of an application sharing type, and the like exist. There is also software which can manage a layout of a plurality of window groups.\nFor example, the network camera has a UI for switching a photographed video image to an image size such as QVGA, VGA, or SXGA and displaying and a UI for providing camera functions such as control dialogue of PTZ, event displaying dialogue, panorama screen control dialogue, and the like to the user (for example, network camera VB-M40 made by Canon Inc).\nFurther, what is called a viewer PTZ or the like in which one rectangle is cut out from a video image displayed on a viewer of the client and is displayed as if it was electronically PTZ-driven has also been realized (for example, network camera VB-M600 made by Canon Inc).\nAs a function of the OS such as Windows (registered trademark) or a function of Explorer, there is such a function that previous window display information is stored and, when the window is displayed again, its display position is restored.\nFurther, there is such a function that by software having a plurality of operating modes, layout information of windows which are displayed on a screen is stored every operating mode, when the operating mode is switched, they are displayed in accordance with the stored layout information (publication of U.S. patent application 2006/0233192). Therefore, each time the display mode is switched, the positions can be restored and the windows can be displayed by the layout which is conditional on the display mode.\nWhen a plurality of windows or dialogues are displayed, they can be also automatically arranged (Japanese Patent Application Laid-Open No. H11-15626).\nIn an image displaying apparatus in the monitoring camera system, display units of monitor video images are important. Further, among them, since the video image is carefully viewed in accordance with a situation or the camera operation is necessary, a watching location and functions to be used are not fixed. Particularly, in the case of enlarging an image size for the purpose of monitoring, a highlighted portion is important.\nHowever, there is such a problem that even if the positions are merely restored to a layout which has previously been conditional on, the images cannot be always displayed in a desirable layout."}
{"text": "Within any enterprise, regardless of the number of memos or instructions relating to document management or retention, there is the possibility that not every document gets stored in a designated storage system. A person may create a document and save it locally on their computing device without saving the document to a designated storage system such as a file server. This may create a problem if the computer having the only copy of a document crashes or the document is otherwise lost.\nDistribution of a document may also lead to problems. A person may distribute a document through e-mail or other communications channel leading to duplicate copies of the document residing across a network. This can be problematic when documents are gathered (e.g., for legal compliance) as multiple copies of the same document may be collected, leading to increased use of storage resources to store the duplicate copies and processing resources to gather and de-duplicate the documents.\nFurthermore, documents are sometimes distributed in a manner that relies on the recipient to handle the document properly with respect to compliance or workflow. For example, a user may forward a document to another user without specifically tying the document to a workflow process or without notifying the recipient that the document should not be forwarded outside of an entity's network. Consequently, it can be difficult for an enterprise to manage processes and ensure that documents are handled properly from the standpoint of compliance, workflow and other processes."}
{"text": "1. Field of the Invention\nThe invention relates to devices with combinations of features that assist people with disabilities in their health monitoring, communication, organization, and navigation needs.\n2. Background\nAs important as electronic devices have become in most people\"\"s daily lives, they can be difficult to use for some major segments of the population. For example, the elderly and persons with mild disabilities. Also persons with little or no experience in the operation of electronic devices often find it difficult to develop the habit of using such devices.\nPersonal objects are, for the purposes of this application, items that a person, by habit or necessity, finds indispensable and, consequently, keeps constantly on their person or within easy reach.\nPersonal objects that are used to improve their lifestyle or cope with a disability, form a base upon which the present invention is built. Examples of personal objects include: canes, hats, tobacco tins, jackets, datebooks, pill boxes, knitting bags, pocketbooks, eyeglass cases, walkers, and orthotic appliances. It is a goal of the invention to provide an electronic assistant that is integrated into the personal object in a manner that augments or improves the standard function of the personal object and insures the electronic assistant is never far away. The functional features of the electronic assistant may include:\nCommunication functions, such as a cell phone, pager, radio transceiver, emergency pager/alert device, etc.;\nHealth monitoring functions, such as heart rate, blood pressure, electrocardio profile, skin conductivity, eye function, chemical or drug concentration, etc.;\nOrganization functions, such as those functions commonly found in a personal digital assistant (PDA); and\nNavigation functions, such as provided by a global positioning system, compass, barometric altimeter, etc.\nThe device\"\"s input interface can be manual, radio frequency, infrared, voice, hard-wire, etc., the output interface can be visual, auditory, etc., and input/output can be either direct or indirect via attached hardware.\nIt is another object of the invention to provide an electronic assistant, comprising; a personal object having at least one electronic module incorporated therein, and each electronic module is selected from the group of consisting of a PDA, a cell phone, a navigation module, a biosensor module, and an emergency alert module.\nIt is another object of the present invention to provide an electronic assistant, comprising; a personal object having at least one electronic module incorporated therein. The personal object is one of a cane, a walking stick, a walker, a wheelchair, a personal transportation vehicle, a purse, a key holder, a watch, a pendant, a hearing aid, an eyeglass frame, and a crutch. Each electronic module is selected from the group of consisting of a PDA, a cell phone, a navigation module, a biosensor module, and an emergency alert module. The PDA includes storage and retrieval of appointments and at least one of addresses, phone numbers, calendar events, memos, and to-dos and other lists. The navigation module includes at least one of a GPS, an altimeter, an electronic compass, and a stored map. The biosensor module includes; at least one biosensor, a biosensor information processor, and at least one biosensor communication channel connecting each biosensor with the biosensor information processor. Each biosensor is one of; a biosensor disposed on a surface of the personal object, a remote biosensor attached to the body of the user, and a remote biosensor surgically implanted in the body of the user. The biosensor communication channel is at least one of; a wire, a radio frequency transmitter and receiver, and an infrared transmitter and receiver. The emergency alert module includes an emergency button, and an emergency notification sequence stored in the electronic assistant capable of calling for assistance through said cell phone.\nIn summary the present invention is an electronic assistant, capable of overcoming a person\"\"s resistance to using electronic devices by incorporating electronic functions in a personal object that the user finds indispensable. Other examples of personal objects include a walking stick, a wheelchair, a personal transportation vehicle, a purse, a key holder, a watch, a pendant, a hearing aid, an eyeglass frame, and a crutch. Some examples of assistant electronics to be integrated with the personal object include a PDA, a cell phone, a navigation module, a biosensor module, and an emergency alert module. The PDA can include storage and retrieval of appointments, phone numbers, calendar events, memos, and to-dos and other lists. The navigation module might include a GPS, an altimeter, an electronic compass, and/or a stored map. Biosensors may be directly incorporated into the personal device or communicate with a receiver therein if the biosensor is required to be attached directly to the user\"\"s body or surgically implanted. The emergency alert module may include an emergency button, and an emergency notification sequence stored in the electronic assistant capable of calling for assistance through the cell phone.\nThe invention will be described in connection with certain preferred embodiments, with reference to the following illustrative figures so that it may be more fully understood. With reference to the figures, it is stressed that the particulars shown are by way of example and for purposes of illustrative discussion of the preferred embodiments of the present invention only, and are presented in the cause of providing what is believed to be the most useful and readily understood description of the principles and conceptual aspects of the invention. In this regard, no attempt is made to show structural details of the invention in more detail than is necessary for a fundamental understanding of the invention, the description taken with the drawings making apparent to those skilled in the art how the several forms of the invention may be embodied in practice."}
{"text": "This invention relates to a device to be used for tending straps and also to be used as a whistle. More particularly, this invention relates to a device for tending straps with a whistle integrated within the device where both the whistling and strap-tending functions are independently operable.\nDevices for tending straps take many forms. The buckle on a typical belt is maybe the most recognizable, but fasteners with similar functions encompass a wide range of designs and related functions, see U.S. Pat. No. 5,794,316 for example. The xe2x80x9ctending strapxe2x80x9d function is intended to encompass the various functions performed by devices that are attached to straps, or more specifically to belts, webbing, or cords, and that work to secure the straps under tension (a buckle or webbing adjuster for example), confine them to a particular area (a slider), or limit movement of other items along the strap (a cord lock). The articles utilizing these designs are ubiquitous where one finds hikers, boaters, campers, climbers, and other persons employing similar gear. Many of these activities give rise to emergencies or other situations where the spoken word either does not carry far enough, is not loud enough to be heard over surrounding noise, or is not distinct enough to draw attention. It is often recommended that the participant carry a whistle for safety in activities where emergency situations can exist. Such needs for help, however, tend to be rather unexpected. This unexpected quality results in people not equipping themselves with whistles at all times and thus not having one when such an emergency occurs. If the whistle could be incorporated into an otherwise useful item then a person would be more likely to have one available at all times, particularly during emergencies. Aside from emergency uses, a whistle may be merely a handy means for the user to notify others of his or her presence, or otherwise get attention. There is thus a need for a device that has a function with a utility related to the activity that also incorporates a whistle that can be used relatively quickly when the need arises.\nWhistles in general take many forms, for examples see U.S. Pat. Nos. 5,546,887 and 5,507,246. Such whistles undeniably perform the necessary function if they are available. But such whistles are often not available for any number of reasons. A solitary whistle is not an item that most would habitually carry. Therefore a potential user would probably need a reason to include a whistle in the gear for that day. Most of the obvious reasons are related to emergencies and most people do not anticipate having emergencies. This would make a whistle a low priority item on anyone\"\"s packing list. And because emergencies are fortunately rare, the whistle would in fact generally be just another item cluttering up a purse, key-chain, or related devices. This limited need makes it easy and understandable to simply forget to carry a whistle even should one think it a good idea in general.\nBut many other things are carried constantly and without additional effort. Many of the activities giving rise to the need for a whistle have a particular type of gear associated with that activity. Although the various activities do not necessarily have gear in common, the individual items of gear have common elements. One group of common elements comprises the fasteners used to secure webbing, belts, straps, and cords on this gear. These devices range from your typical belt buckles to the more high tech molded side release buckles, tension adjusters, sliders, and cord locks. They are present in gear ranging from book bags and backpacks, to key-chains, bike helmets, and the cords on the hoods of many jackets.\nBuckles have been proposed incorporating tools such as a whistle, however the design has been unsatisfactory and has not gained wide spread acceptance. For example, U.S. Pat. Nos. 3,885,250 and 3,903,547 disclose a buckle with a prong that may be fashioned into a whistle. While specific details of how such a whistle might be formed are not disclosed, the general arrangement is not satisfactory because, being on a prong, such a whistle would be hidden and even if known, use would require disengagement of the retained strap or belt. Additionally, the strap must be large enough to house the whistle. This requirement prevents incorporation of a whistle where the strap is of small cross-section because the prong must be accommodated within either a thick strap or two layers of strap. This also means that the strap must be specially manufactured to accept the flange, increasing manufacturing costs and reducing the application of the buckle because it must necessarily be part of a matched set to function. Finally, since the prong is rigid and extends into the strap the areas where the buckle can be located are limited to relatively flat locations that extend the length of the prong and buckle.\nThus there exists a present need for a device with an integral whistle that can function with a variety of strap sizes and shapes to fasten, tension, lock or generally tend the straps and where the whistle is unobtrusive yet apparent and accessible to the user.\nThe preferred embodiments of the present invention combine the function of a whistle with that of a device for tending straps, cords, belts, or webbing of various shapes and sizes. The whistle is combined with the strap-tending device in a manner that allows the whistle to be used even when the device is tending the straps as designed. Given that a strap runs in a general direction and that the strap-tending device typically performs a function that is oriented along the same general direction, the placement of the whistle is preferably transverse to that direction to facilitate access to the whistle, decrease device dimensions, and improve its manufacturability. The strap-tending functions of the preferred embodiments include: a two-piece side release buckle with the whistle incorporated in the female half and the strap adjusting mechanism in the male half; a single piece strap adjuster with the whistle located opposite the anchor bars and slots from the extension designed to facilitate releasing tension; a cord lock with the whistle integrated into the barrel; and a slider with the whistle integrated to one side of the slots designed to engage the straps. In all cases the whistle is accessible and functional with the intended strap or straps in place and the device performing its intended strap-tending function."}
{"text": "The present invention relates to a coin acceptance means and method, and, more particularly, to a coin acceptance means for use in a vending system wherein a vend price is established prior to coin deposit or the submission of other credit means by a customer, for determining, based upon the ability of the system to provide a correct amount of change to the customer if a particular deposited coin or other credit submission is accepted, whether a given deposited coin or credit submission, even if valid, will be accepted.\nIt will be appreciated that, throughout this application, the term \"coin\" may be employed to mean any coin (whether valid or counterfeit), token, slug, washer, or other item, including currency and scrip, which might be used by an individual in an attempt to operate a coin-operated device or system. A \"valid coin\" is considered to be an authentic coin, token, or the like of the monetary system or systems in which or with which the coin-operated device or system is intended to operate and of a denomination which the device or system is intended selectively to receive and to treat as an item of value. A coin is considered to be \"accepted\" when it is determined to be a valid coin and a credit therefor is entered in the vending system.\nGenerally, single price vending systems have a pre-established vend price against which an accumulated credit is compared to determine whether a vend operation will be effected. Many multi-price vending systems, on the other hand, are so designed to permit the customer to deposit coins to build up an accumulated credit entry, which credit entry is compared against a vend price subsequently established when the customer actuates a vend selection means to select a particular product having a particular vend price. However, even some multi-price vending systems may be so designed that a vend price is established prior to any coin deposit or any recognition thereof. For example, a multi-price system may be so designed to require selection of the desired product prior to the recognition of any coin deposit. With such a system the vend price would be established by the product selection made, and the vending system would thereafter function in similar fashion to a single price vending system. The coin acceptance means of the present invention is primarily directed to single price vending systems and to multi-price vending systems wherein a vend price for a given vend operation is established prior to the deposit of coins or submission of credit with regard to such particular vend operation. For the sake of convenience, such systems will hereinafter be referred to as predefined price vending systems.\nIn the past, most predefined price vending systems have functioned in one of two different ways to try to ensure that a customer would receive the proper amount of change in the event that he had made an overdeposit. A number have been designed such that an exact change light is caused to be illuminated whenever the change making tube for the lowest denomination coin is found to be empty. With some of such vending systems, especially systems wherein the vend price is equal to the value of a particular denomination of coin, if the customer chooses to ignore the exact change light and to deposit a coin which does not constitute \"exact change\", the system will refuse to accept such coin and will simply return such deposited coin to the customer. In such a way, such systems can ensure that no overdeposit will be accepted and that no payback will be required. With other \"exact change\" systems, if the customer chooses to ignore the exact change light and deposits a coin sufficient to result in an overdeposit situation, the system will function to accept the coin and to effect a vend of the product selected, and will also attempt to provide a refund in the proper amount, but with no guarantee that the correct amount, or, perhaps, any amount, will actually be refunded to the customer. Consequently, in such a situation, it is quite likely that the customer may receive less than the full refund amount to which he is really entitled.\nOther vending systems have been so designed that they receive the deposited coin and total the credit entries, subtract the vend price from such credit total to obtain a refund amount due, and then determine if sufficient change remains available in the change making tubes to actually pay back to the customer the refund amount due. Such types of vending systems have generally allowed the customer to deposit coins prior to the making of any vend selection in order to accumulate a credit entry, and have permitted such customer to thereafter make a vend selection to establish a vend price. Such systems typically determine both whether the accumulated credit entry is sufficient in view of the vend price established and whether any refund amount will be due to the customer as a result of an overdeposit. If it is determined that the credit entry is at least equal to the vend price and that sufficient change remains to pay back any refund amount due, such systems have then collected the deposited coins, effected a vending of the selected product, and paid back the particular refund amount, if any, due. However, if they have determined that insufficient change remains to pay back the particular refund amount that would be due, all the deposited coins, or coins of like total value, have been returned to the customer and no product vend has been effected.\nRepresentative of some of the various types of vending systems noted hereinabove and of controls therefor or coin changing features that could be employed therewith are those constructions disclosed in U.S. Pat. Nos. 4,188,961; 4,191,999; 4,462,512; 4,499,982; and 4,499,985.\nWith many of the known predefined price vending systems that include changemakers, a vend is not permitted if proper coins are not available for use in refunding the amount of an overdeposit. In such cases, the customer, who has typically deposited a plurality of coins to accumulate a credit entry, which coins are retained by the system upon deposit and the validation thereof, is refunded the total amount of his deposit, either automatically, as a function of the vending system, or upon actuation by him of an escrow switch. In any event, the entire deposited amount of credit is returned to the customer. If, by chance, such customer has available to him additional or other change in such an amount that he can deposit an exact change amount or some other amount for which the vend system includes appropriate change to allow a refund to be made, such customer may seek to obtain a vend by again commencing coin deposit, often re-depositing many of the same coins just returned to him. The noted return of all of such coins to the customer and the subsequent re-deposit of a portion thereof by him in further efforts to obtain a vend may take a significant amount of time in terms of the ability of the vend system to vend products to customers. Such time factor may become problemsome if there are a number of people waiting to use the vend system and several of the people near the beginning of the line all initially make deposits for which the system cannot provide a correct refund. If even some of such customers, after a first refund of all deposited coins, seek to re-deposit coins to obtain a vend, the delay to customers farther back in line will be increased. If the delay becomes too long, the individuals farther back in line may grow irritated, and, in some instances, decide not to wait for their turn, as a consequence of which sales will be lost. Even those who are willing to wait for their turn may develop less positive attitudes toward the vend system and the products vended thereby, the long term effect of which may be a reduction in sales. It is therefore desirable that such delay time in vending, which delay time results from the return of all deposited coins when non-refundable overdeposits are detected and the re-deposit of a portion of such returned coins in a subsequent vend operation, be minimized or eliminated, and the present invention is directed to such end."}
{"text": "Generally, in a plant, various equipment are arranged including electric equipment and monitor control equipment. In devices including such equipment, many require long-term stable operation. It is thus useful to prevent operational stoppages caused by failure of the equipment.\nTo prevent such stoppages due to equipment failure, it is possible to record a failure history of a device, and preventive maintenance can be performed as necessary based on the tendency of a failure to occur.\nAs a product and related information management method, Japanese Patent Disclosure (kokai) No. 2003-308363 discloses a method in which a guarantee on a product or data relating to the product, such as management information of a recycle cost, can be easily read.\nNamely, a product includes a memory unit for readably storing data relating to the product. The data relating to the product includes product name and/or the manufacture's serial number, and the term of guarantee on the product corresponding to the product name and/or the manufacture's serial number.\nIn the related information managing method, the information relating to the product is managed by means of a computer, an identifier for identifying each of the products is readably stored in a memory unit in the product, the identifier is read, and information relating to the corresponding product is searched with this identifier as a search key and displayed.\nIn Japanese Patent Disclosure (kokai) No. 2003-308363, an identifier for identifying each of the products is readably stored in a memory unit of the product, the identifier is read, and information relating to the corresponding product is searched with this identifier as a search key and displayed.\nTherefore, in Japanese Patent Disclosure (kokai) No. 2003-308363, failure history data of the shipped product is not managed.\nParticularly, in a large-scale plant, intricately and highly developed devices are arranged. It would be useful for an equipment manager to manage failure history data of these devices, perform preventive maintenance as necessary based on a tendency of failures to occur, and prevent the stoppage of the equipment due to the failure ahead of time."}
{"text": "Advances in semiconductor manufacturing technology have resulted in, among other things, reducing the cost of sophisticated electronic products to the extent that the integrated circuits fabricated in accordance with these semiconductor manufacturing technologies have become ubiquitous in the modern environment.\nAs is well-known, integrated circuits are typically manufactured in batches, and these batches usually contain a plurality of semiconductor wafers within and upon which integrated circuits are formed through a variety of semiconductor manufacturing steps, including, for example, depositing, masking, patterning, implanting, etching, planarizing and so on.\nCompleted wafers are tested to determine which die, or integrated circuits, on the wafer are capable of operating according to predetermined specifications. In this way, integrated circuits that cannot perform as desired are not packaged, or otherwise incorporated into finished products.\nIt is common to manufacture integrated circuits on roughly circular semiconductor substrates, or wafers. Further, it is common to form such integrated circuits so that conductive regions disposed on, or close to, the uppermost layers of the integrated circuits are available to act as terminals for connection to various electrical elements disposed in, or on, the lower layers of those integrated circuits. Such conductive regions are commonly referred to as pads. Pads are commonly used to provide electrical access to the integrated circuit both during testing and during the operation of the integrated circuit as it is incorporated into a final product. Typical pads are formed from aluminum. It is well-known that the surfaces of aluminum which are exposed to the atmosphere will oxidize. These oxidation layers interfere with the formation of low resistance electrical connection to the pads.\nDuring the process of testing the performance of the integrated circuits, the pads are commonly contacted with probe needles, or other contact structures, of a probe card. It is through such temporary connections that a test apparatus may electrically interact with an integrated circuit.\nConventionally, the probe needles, or other contact structures, of the probe card are brought into physical contact with the pads and then moved laterally to “scrub” the pad. The scrubbing operation is intended to break through the oxide layer on the top surface of the pad, thereby providing reduced contact resistance. Unfortunately, scrubbing disturbs the pad structure and can contribute to yield loss due to failure of bond wires to properly attach to the disturbed pad structure.\nWhat is needed are methods and apparatus for providing low contact resistance connections to pads of integrated circuits without disturbing or substantially redistributing the material from which the pads are formed."}
{"text": "The present invention relates in general to hopper transport vehicles, and, more particularly, to a bottom for the hoppers of such vehicles.\nHopper vehicles include trucks, railroad cars, and the like, and are used to transport many materials, such as granular products, liquid products, and the like. The problems involved in unloading such vehicles have engendered many inventions, see, for example, co-pending Patent Application Ser. Nos. 137,438, filed Apr. 3, 1980, and 144,850, filed Apr. 29, 1980 by the present assignee, and incorporated herein by reference thereto.\nMany unloading systems involve pressurizing the hopper and forcing the product into a discharge manifold. Such unloading systems will be referred to hereinafter as pneumatic unloading systems. Other unloading systems involve gravity unloading wherein product is discharged out of the hopper bottom.\nThere are situations where pneumatic unloading is more efficient than gravity unloading, and vice versa. Accordingly, there is need for a hopper vehicle which is capable of being unloaded by pneumatic unloading procedures and by gravity unloading procedures."}
{"text": "Bioassays are used to probe for the presence and/or the quantity of a target material in a biological sample. In surface based assays, the target amount is quantified by capturing it on a solid support and then detecting it. One example of a surface-based assay is a DNA microarray. The use of DNA microarrays has become widely adopted in the study of gene expression and genotyping due to the ability to monitor large numbers of genes simultaneously (Schena et al., Science 270:467-470 (1995); Pollack et al., Nat. Genet. 23:41-46 (1999)). More than 100,000 different probe sequences can be bound to distinct spatial locations across the microarray surface, each spot corresponding to a single gene (Schena et al., Tibtech 16:301-306 (1998)). When a fluorescent-labeled DNA target sample is placed over the surface of the array, individual DNA strands hybridize to complementary strands within each array spot. The level of fluorescence detected quantifies the number of copies bound to the array surface and therefore the relative presence of each gene, while the location of each spot determines the gene identity. Using arrays, it is theoretically possible to simultaneously monitor the expression of all genes in the human genome. This is an extremely powerful technique, with applications spanning all areas of genetics (For some examples, see the Chipping Forecast supplement to Nature Genetics 21 (1999)). Arrays can also be fabricated using other binding moieties such as antibodies, proteins, haptens or aptamers, in order to facilitate a wide variety of bioassays in array format.\nOther surface-based assays include microtitre plate-based ELISAs in which the bottom of each well is coated with a different antibody. A protein sample is then added to each well along with a fluorescent-labeled secondary antibody for each protein. Target proteins are captured on the surface of each well and secondarily labeled with a fluorophore. The fluorescence intensity at the bottom of each well is used to quantify the amount of each target molecule in the sample. Similarly, antibodies or DNA can be bound to a microsphere such as a polymer bead and assayed as described above. Once again, each of these assay formats is amenable for use with a plurality of binding moieties as described for arrays.\nOther bioassays are of use in the fields of proteomics, and the like. For example, cell function, both normal and pathologic, depends, in part, on the genes expressed by the cell (i.e., gene function). Gene expression has both qualitative and quantitative aspects. That is, cells may differ both in terms of the particular genes expressed and in terms of the relative level of expression of the same gene. Differential gene expression is manifested, for example, by differences in the expression of proteins encoded by the gene, or in post-translational modifications of expressed proteins. For example, proteins can be decorated with carbohydrates or phosphate groups, or they can be processed through peptide cleavage. Thus, at the biochemical level, a cell represents a complex mixture of organic biomolecules.\nOne goal of functional genomics (“proteomics”) is the identification and characterization of organic biomolecules that are differentially expressed between cell types. By comparing expression, one can identify molecules that may be responsible for a particular pathologic activity of a cell. For example, identifying a protein that is expressed in cancer cells but not in normal cells is useful for diagnosis and, ultimately, for drug discovery and treatment of the pathology. Upon completion of the Human Genome Project, all the human genes will have been cloned, sequenced and organized in databases. In this “post-genome” world, the ability to identify differentially expressed proteins will lead, in turn, to the identification of the genes that encode them. Thus, the power of genetics can be brought to bear on problems of cell function.\nDifferential chemical analyses of gene expression and function require tools that can resolve the complex mixture of molecules in a cell, quantify them and identify them, even when present in trace amounts. The current tools of analytical chemistry for this purpose are presently limited in each of these areas. One popular biomolecular separation method is gel electrophoresis. Frequently, a first separation of proteins by isoelectric focusing in a gel is coupled with a second separation by sodium dodecyl sulfate-polyacrylamide gel electrophoresis (SDS-PAGE). The result is a map that resolves proteins according to the dimensions of isoelectric point (net charge) and size (i.e., mass). Although useful, this method is limited in several ways. First, the method provides information only about two characteristics of a biomolecule-mass and isoelectric point (“pI”). Second, the resolution power in each of the dimensions is limited by the resolving power of the gel. For example, molecules whose mass differ by less than about 5% or less than about 0.5 pI are often difficult to resolve. Third, gels have limited loading capacity, and thus limited sensitivity; one often cannot detect biomolecules that are expressed in small quantities. Fourth, small proteins and peptides with a molecular mass below about 10-20 kDa are not observed.\nThe use of functionalized chips is replacing gels as the method of choice for bioassays. Efforts to improved the sensitivity of assays have resulted in a number of chip designs. For example, a specific binding assay device, which comprises multilayer analytical materials is known (see, for example, EP 51183, EP 66648, DE 3227474 and EP 236768). other multilayer chips are set forth in U.S. Pat. Nos. 4,839,278 and 4,356,149.\nAn effective chip for bioassay applications must have adequate capacity to immobilize a sufficient amount of an analyte from relevant samples in order to provide a suitable signal when subjected to detection (e.g., mass spectroscopy analysis). Suitable chips must also provide a highly reproducible surface in order to be gainfully applied to profiling experiments, particularly in assay formats in which the sample and the control must be analyzed on separate adjacent chip surfaces. Chips that are not based on a highly reproducible surface chemistry result in significant errors when undertaking assays (e.g., profiling comparisons).\nIn general, currently available chips are based on in-situ polymerization of a hydrogel on each spot of a bioassay chip (see, e.g., WO 00/66265 (Rich et al.) The selectivity and reproducibility of these chips is highly dependent upon a large number of experimental variables including, monomer concentration, monomer ratios, initiator concentration, solvent evaporation rate, ambient humidity (in the case when the solvent is water), crosslinker concentration, laboratory temperature, pipetting time, sparging conditions, reaction temperature (in the case of thermal polymerizations), reaction humidity, uniformity of ultraviolet radiation (in the case of UV photopolymerization) and ambient oxygen conditions. While many of these parameters can be controlled in a manufacturing setting, is difficult if not impossible to control all of these parameters which impinge upon reproducibility. As a result, in situ polymerizations results in relatively poor reproducibility of all parameters including spot to spot, chip to chip and lot to lot.\nThus, there is a tremendous need for chips that provide reproducible results from assay to assay, are easy to use, and provide quantitative data in multi-analyte systems. Moreover, to become widely accepted, the chips should be inexpensive to make, and to use for the detection of analytes. The availability of a chip having the above-described characteristics would significantly affect research, individual point of care situations (doctor's office, emergency room, out in the field, etc.), and high throughput testing applications. The present invention provides chips having these and other desirable characteristics"}
{"text": "A data center may house computer systems and associated components, such as telecommunications systems, servers, and storage systems. A data center may also include redundant or backup components, such as power supplies, data communications connections, environmental controls, security devices, and the like. A data center may operate systems associated with a particular company. For example, a company may hire a third-party service provider to manage the company's systems through the third-party service provider's data center."}
{"text": "The present invention relates to systems to attenuate a selected electrical signal in response to changes in selected conditions including changes in characteristics of the signal. Signal attenuation means are utilized in numerous applications from highly sophisticated applications such as computers to more straightforeward applications in simple electrical or electromechanical devices.\nSignal attenuation can, of course, be accomplished manually by observation of the situation to be controlled and manual manipulation of the signal control means. In the alternative, signal attenuation can be accomplished by sensing devices responsive to the condition to be controlled and use of controlled devices to automatically attenuate signals in response to changed conditions.\nIn one application of the foregoing, specifically in electronic circuitry, attenuation can be accomplished by the use of amplifier means for the controlled signal with continuous feedback and gain adjustment of the amplifier circuit. Such arrangements provide satisfactory control in some instances but in other application, for example in attenuation of the output of multiple input sound actuated sound amplification systems, such arrangements have been less satisfactory than desired. Also, such systems are subject to distortion resulting from transient and feedback effects.\nOther attenuation methods have provided for stepwise, rather than continuous attenuation. Such arrangements have generally been manually controlled and where automatic control has been provided as disclosed in an article entitled \"Automatic Attenuator Rapidly Changes Signal Level\", Electronics, Sept. 15, 1969 at page 120 by continuous monitoring of incoming signals to shift attenuation levels. Such systems have been found less than desirable in sound amplification because the rapid switching of output channels leads to undesirable interference with the quality of transmitted program material.\nIn some sound attenuation devices so called voltage controlled amplifiers have been utilized wherein the attenuation or amplification of a program signal is accomplished by varying the voltage to the gain control of the amplifier to modify the attenuator to maintain selected output signal characteristics. However, such voltage controlled amplifiers have inherently poor noise factors and tend to distort the program to be amplified."}
{"text": "The invention relates to an instrument for dental use, of the type comprising a handle and one or two active ends that are pointed or variously shaped for use in various types of procedures, and in particular to a specific type of finish on the handle and a method of forming the finish.\nCurrently known instruments of this type have a handle of which a portion has a knurled or ridged surface, in particular a knurled surface with deep indentations. The portion of the surface machined in this way is currently thought to be necessary to enable the dentist to have a firm grip on the instrument.\nThe knurling, while considered necessary for gripping the instrument, has proved to be a source of major disadvantages, particularly with regard to hygiene. The indentations in the knurled surface trap contaminants and dirt that could transmit microorganisms. Although they are sterilized, instruments of this type may not be perfectly cleaned because contaminants can accumulate in the indentations in the knurled surface.\nThe primary object of the present invention is a surgical instrument, particularly intended for dental use and method of making same, which does not have the above disadvantages but nevertheless still has the advantage of enabling the dentist to grip the instrument correctly and firmly.\nThe instrument according to the present invention is substantially characterized in that the handle is substantially smooth and has at least one grip zone which has a plurality of microchannels extending in a circumferential direction of the handle. The plurality of microchannels form a uniform surface roughness with the microchannels being arranged on the handle in a uniform random distribution. Each of the micro-channels are sized to be below a magnitude to collect contaminants. The grip zone is formed by providing a grinding surface with a magnitude of roughness that is below a roughness to cause indentations in the handle that would collect contaminants. A grinding surface in the form sandpaper having a grade or sieve number between n. 105 and 500 of the US standard sieves scale is used to create the micro-channels. Sandpaper of n. 240 grade has been found to the most desirable. Other grinding surfaces and other devices can also be used which form the same type of uniform random distribution of microchannels as formed by sandpaper having a grade between n. 105 and n. 500 of the US standard sieves scale. The functional element of the instrument is connected to one of the ends of the handle for performing a dental procedure. The grinding surface is rotated relative to an outer surface of the handle and the handle is held against the grinding surface to form micro-channels in the handle. The relative rotating has a directional component in a circumferential direction of the handle. The micro-channels thus formed are of a size being below a magnitude to collect contaminants. The relative rotating forms the micro-channels in the circumferential direction of the handle.\nAlternatively, the outer surface of the handle is also moved against the grinding surface in an axial direction of the handle to also form the micro-channels in the handle. The axial moving is performed substantially simultaneously with the relative rotating to form the micro-channels in both the circumferential direction and the axial direction of the handle. This causes the micro-channels to be formed in a helicoidal direction of the handle.\nThis type of grip surface provides sufficient roughness to ensure that the user can grip it correctly, and that there is no need for indentations, cavities or other forms of uneven surface that could trap contaminants or pathogens. In the context of the present description the grinding surface covers any surface or device which forms a uniform random distribution of microchannels as formed by sandpaper having a grade between n. 104 and n. 500 of the US standard sieves scale. This size and distribution has been found to allow sufficient roughness to be obtained without the formation of indentations which could form a trap for contaminants or pathogens.\nThe various features of novelty which characterize the invention are pointed out with particularity in the claims annexed to and forming a part of this disclosure. For a better understanding of the invention, its operating advantages and specific objects attained by its uses, reference is made to the accompanying drawings and descriptive matter in which preferred embodiments of the invention are illustrated."}
{"text": "A continuing goal of integrated circuit fabrication is to decrease the dimensions thereof. Integrated circuit dimensions can be decreased by reducing the dimensions and spacing of the constituent features or structures. For example, by decreasing the dimensions and spacing of semiconductor features (e.g., storage capacitors, access transistors, access lines) of a memory device, the overall dimensions of the memory device may be decreased while maintaining or increasing the storage capacity of the memory device.\nAs the dimensions and spacing of semiconductor device features become smaller, conventional lithographic processes become increasingly more difficult and expensive to conduct. Therefore, significant challenges are encountered in the fabrication of nanostructures, particularly structures having a feature dimension (e.g., critical dimension) less than a resolution limit of conventional photolithography techniques (currently about 50 nm). It is possible to fabricate semiconductor structures of such feature dimensions using a conventional lithographic process, such as shadow mask lithography and e-beam lithography. However, use of such processes is limited because the exposure tools are extremely expensive or extremely slow and, further, may not be amenable to formation of structures having dimensions of less than 50 nm.\nThe development of new processes, as well as materials useful in such processes, is of increasing importance to make the fabrication of small-scale devices easier, less expensive, and more versatile. One example of a method of fabricating small-scale devices that addresses some of the drawbacks of conventional lithographic techniques is self-assembled block copolymer lithography.\nAlthough self-assembled block copolymer lithography is useful for fabrication of semiconductor structures having dimensions of less than 50 nm, there are still problems that must be addressed. Self-assembled block copolymer materials may not provide nanostructures exhibiting sufficiently low defect levels.\nSelf-assembled nucleic acids have been researched for forming semiconductor devices. The specificity of complementary base pairing in nucleic acids provides self-assembled nucleic acids that may be used for self-assembled nucleic acid lithography processes."}
{"text": "FIG. 1 is a cross-section of a known formation testing/sampling probe assembly 100 viewed in a plane containing the wellbore axis. The assembly 100 includes a packer 101 configured to be pressed against a formation of interest via hydraulically actuated pistons schematically shown as 103 and 104. Deformation of the packer 101 around an extendable probe barrel 102 and against the formation creates a seal that isolates the formation fluids and pressure from the wellbore environment. The probe barrel 102 extends forward, due to high pressure hydraulic oil entering chamber 105 at the same time pistons 103 and 104 are actuated to apply an extension force. The area of investigation by the probe assembly 100 is limited to the area that is in direct contact with the extendable probe barrel 102, usually called the orifice. This area is limited by the diameter of the probe barrel 102, which is about the same dimension as the diameter of the hole 106 in the packer 101 through which the probe barrel 102 passes.\nOther prior art has sought to improve upon the design of the probe assembly 100 by increasing the diameter of the probe barrel. Nonetheless, the area of investigation is still limited to the largest diameter of the probe barrel, which is currently between two and three inches. This could be problematic for formations with thin laminations where the formation may have a small (e.g., ˜0.5 inches thick) production zone sandwiched between thicker zones of impermeable formation. With prior art designs, finding a non-producing zone is much more likely due to the small lateral extent (the probe barrel diameter) of the area of investigation.\nOne attempt at addressing this thin lamination problem involved elongating the packer, such as with the known probe assembly 200 shown in FIG. 2. The elongated packer 201 of the probe assembly 200 is pressed against the formation by backing plate 202 to create a seal. However, in this case, an annular shaped metal spacer 207a fills the upper hydraulic chamber 207, preventing the probe barrel 208 from extending forward. Since the packer 201 and backing plate 202 are elongated along the central axis of the wellbore, a number of formation laminations could be investigated. The axial length of investigation using the probe assembly 200 could approach several inches, including one known embodiment in which the length was about seven inches.\nHowever, problems with the probe assembly 200 have been discovered. For example, the packer 201 is not as constrained at its inner boundary as it is near the raised rim 203 of the backing plate 202 that defines the orifice of the probe. It has been found that the packer material 201 migrates into the probe's orifice during formation testing operations. The large difference in pressures between the wellbore, pushing on surfaces 204, and the probe orifice, pushing on the packer surface near the raised rim 203, forces the packer material to move into the orifice. This movement decreases the ability of the packer 201 to seal against the formation and prevents the tool's pressure gauges from equilibrating to that of the pressure of the formation fluid. As the elastomeric material of the packer 201 is sucked into the probe orifice, the hydraulic pressure declines in the pistons (schematically shown in FIG. 2 as 205 and 206) applying the extension force (a consequence of the tool design). This, in turn, results in a decline in the force pushing the packer against the formation, further exacerbating the situation, and resulting in both packer damage and loss of seal. This design has lead to a high number of lost seals.\nAnother problem that exists with the design shown in FIG. 2 is that the upper surface of the raised rim 203 of the backing plate 202 is located well below the surface of the elastomer 201 that seals against the wellbore wall. As a consequence, elastomer must be compressed before the raised rim 203 makes contact with the wellbore wall. This is of considerable disadvantage because it requires performing large volume pretests. On the other hand, if this difference in heights between the surface of the raised rim 203 and the sealing surface of the elastomer 201 didn't exist, then it would not be possible for the elastomer 201 to form a seal against the wellbore wall, because the metal spacer 207a makes the probe barrel 208 and the backing plate 202 move in unison."}
{"text": "(1) Field of the Invention\nThis invention relates to a novel method for processing a picture to reproduce the picture with improved visibility on a color television receiver, and to a picture processor therefor.\n(2) Description of the Related Art\nReference is first made to FIG. 3 which is a simplified block diagram of a transmission system for color television pictures. In a television camera 21, a multicolored light is received from an object through a lens 22 and the color of the object is separated into the three primary colors, i.e., red (R) light, green (G) light, and blue (B) light by use of a dichroic mirror unit 23. From these colors of light, primary color signals E.sub.R,E.sub.G,E.sub.B are then produced by image orthicon camera tubes 21R,21G,21B, respectively. A luminance signal E.sub.Y, representative of the brightness of the picture, and color signals having information as to hue and saturation values are reproduced from the primary color signals by a matrix circuit 24 are then transmitted. The color signals which provide the primary color signals E.sub.R,E.sub.G,E.sub.B when combined with the luminance signal E.sub.Y are color difference signals E.sub.R -E.sub.y,E.sub.G -E.sub.Y, E.sub.B -E.sub.Y. When such color difference signals are transmitted, the luminance signal E.sub.Y is added in the following manner at the reception side, thereby making it possible to reproduce the original primary colors. EQU (E.sub.R -E.sub.Y)+E.sub.Y =E.sub.R EQU (E.sub.G -E.sub.Y)+E.sub.Y =E.sub.G EQU (E.sub.B -E.sub.Y)+E.sub.Y =E.sub.B\nAs depicted in FIG. 4, the signals thus transmitted are received by an antenna 25. The desired channel is selected by a tuning circuit 26, so that the signals are amplified and converted into intermediate frequency signals. Video intermediate frequency (IF) signals which have been amplified at an intermediate frequency (IF) amplifier 27 are converted into video signals at a video signal detector 28. After amplification of the video signals at a video signal amplifier 29, the resulting signals are fed to a color picture tube 1. By electron beam scanning, a complete picture is formed from a series of video signals. Here, synchronizing signals transmitted along with the video signals are separated at a synchronous circuit 31 to control a horizontal and vertical deflection circuit 32, whereby the frequency and phase of the electron beam scanning are maintained consistent with those at the transmission side.\nOn the other hand, audio intermediate frequency (IF) signals are converted into audio FM signals at an audio signal detector and FM converter 33. The audio FM signals are converted further into audio signals by an FM detector 34 and a low frequency (LF) amplifier 35 and are then reproduced at a speaker 36.\nThe video signal amplifier 29 serves to amplify each luminance signal of the color television signals, which have been obtained at the video signal detector 28, to a level sufficient to activate the color picture tube 1. The video signal amplifier 29 is usually constructed of amplification circuits arranged in 3 to 5 stages and various auxiliary circuits.\nTo reproduce a color picture, the color picture tube 1 has to be fed with primary color signals. There are two types of manners for feeding these signals, i.e., the primary color drive system and the color difference drive system. FIG. 5 is a block diagram showing how primary color signals are produced and then fed to the color picture tube 1 in the primary color drive system. Primary color signals E.sub.R,E.sub.G,E.sub.B are applied at a negative polarity to respective cathodes 37-1,37-2,37-3 of the color picture tube 1, while respective first grids 38-1,38-2,38-3 are fed with a d.c. voltage.\nThe primary color signals are produced by adding the luminance signal and color difference signals at a matrix circuit 39. The matrix circuit 39 is inputted with the luminance signal E.sub.Y and color difference signals E.sub.R -E.sub.Y,E.sub.G -E.sub.Y,E.sub.B -E.sub.Y and produces as outputs negative primary color signals -E.sub.R,-E.sub.G,-E.sub.B. The negative primary color signals are then fed to the respective cathodes 37-1,37-2,37-3 of the color picture tube 1.\nIn the color picture tube 1, electron beams controlled by the primary color signals are caused to impinge exactly on color phosphors of red, green and blue colors, respectively, by way of a color sorting system, free of illustration, so that the phosphors are caused to glow to reproduce the color picture. The construction of the color picture tube 1 is illustrated in FIG. 6. The color picture tube 1 is basically constructed of three components, namely, an electron gun 61 for generating and directing electron beams, a shadow mask 62 for causing the electron beams to impinge only on phosphor dots of the designated colors, and a phosphor screen 63 coated precisely with red (R), green (G) and blue (B) phosphor dots in a predetermined pattern. Depending on differences in these components, color picture tubes can be classified into several types. Their details are however omitted herein.\nFormation of a picture on a color picture tube is performed by left-to-right scanning (horizontal scanning) and up-to-down scanning (vertical scanning). The picture is successively and downwardly formed on the picture screen from the upper left corner. When the scanning has reached the lower right corner, the scanning returns again to the upper left corner. In this manner, 30 pictures, which are each divided into 525 lines, are formed per second. Accordingly, the horizontal scanning is repeated 15,750 times a second (525 lines.times.30 pictures) while the vertical scanning is repeated 30 times.\nHowever, interlaced scanning is used in actual scanning. As a result, the vertical scanning is repeated 60 times. As illustrated in FIG. 7(a) and 7(b), interlaced scanning is performed by first coarsely scanning at some intervals as indicated by numbers 1, 2, 3 and 4 and then scanning once more between the first scanning lines as indicated by numbers 5, 6 and 7. The scanning of an entire picture is brought to completion by only these two vertical scanning operations. When interlaced scanning is performed, two coarse pictures of 262.5 scanning lines are sent successively to form a picture of 525 scanning lines. It therefore appears as if 60 pictures are sent out, although 30 pictures are actually sent out per second. It is hence possible to obtain pictures of little flicker.\nAs has been described above, the conventional television receivers have reached a technical level which is satisfactory from the practical viewpoint. It is, however, now strived to improve the dissolution further, for example, by increasing the number of scanning lines to 650 or 700 lines per screen or converting the signal processing circuits into their digital counterparts. In addition, high-definition television receivers have also been developed, as future television receivers, for commercial use."}
{"text": "1. Field of the Invention\nThis invention relates to cementing plugs for use in cementing casing in a well, and more particularly, to a universal cementing plug having improved wiping and extended wear and which includes a plurality of interchangeable inserts so that the plug may be selectively used as a top or a bottom plug.\n2. Description of the Prior Art\nIn the process of preparing a well for testing and/or production, a casing is positioned in the well and cemented in place. Typically, at the beginning of the cementing job in rotary-drilled wells, the casing and the wellbore are usually filled with drilling mud. In many areas, to reduce contamination on the interface between the mud and cement a bottom plug is released from a plug container and pumped ahead of the cement slurry. Such plugs have wipers of an elastomeric material thereon to wipe the casing of any accumulated mud film so that the mud is pushed ahead of the bottom plug.\nWhen the bottom plug reaches floating equipment such as a float collar or float shoe at the bottom of the casing string, a fluid pressure differential created across the plug ruptures a rubber diaphragm at the top of the plug and allows the cement slurry to proceed down the casing through the plug and floating equipment and then up an annulus space defined between the casing and the wellbore.\nWhen all of the cement has been mixed and pumped into the casing string, a top cementing plug is released from the plug container. The top plug also has wipers of elastomeric material thereon. The function of the top plug is to follow the cement and wipe any accumulated cement film from the inside of the casing. The top plug is also designed to reduce the possibility of any contamination or channeling of the cement slurry with the drilling mud that is used to displace the cement column down the casing and into the annular space between the casing the wellbore. The top cementing plug is typically solid in construction, and the design is such that when it reaches the bottom cementing plug at the float collar or float shoe, the top cementing plug causes a shutoff of fluids being pumped into the casing. This causes a normal pressure rise at the surface and notifies the operator that the cementing job is complete.\nThe landing of the top plug lessens the possibility of any further displacement of the cement slurry and provides a better quality of cement slurry around the bottom of the casing where a good cement bond to the casing is required.\nCurrently, two different cementing plugs are used in this cementing operation, one for the top and one for the bottom. The bottom plug has a shearable member, such as the rubber diaphragm previously mentioned, which shears when a specific fluid pressure differential is applied thereto. The top plug is substantially solid. Because each plug requires different construction, separate molds must be used for each of the plugs which increase the costs of manufacturing, and also, the two separate plugs must be kept in inventory. The present invention solves this problem by using a single plug subassembly design which has the same general construction whether it is used as a top plug or a bottom plug. A shearable insert in positioned in one plug so that it may be used as a bottom plug. This shearable member is designed to shear at a predetermined differential pressure thereacross. In one embodiment, the shearable member is a flat disc, and in another embodiment, the shearable member has a relatively thin domed portion. Another insert, which is essentially non-shearable at the pressures in which the plugs are utilized, is positioned in another plug so that it can be used as a top plug. By the use of a single plug subassembly, with separate inserts, the cost of molds of the plugs is decreased, and only one plug must be maintained in inventory along with the different inserts.\nAnother advantage of the present invention is that the shearable member may be interchangeable with a plurality of shearable members, including, but not limited to, the two embodiments previously described, designed to shear at any one of a selected number of differential pressures as necessary for different well conditions. This is an improvement over the previous design which had essentially one shear pressure.\nWith prior art cementing plugs, the wiping efficiency of the wipers on the plugs is affected by pumping rate and wear along the casing surface. The cementing plug of the present invention provides an improved wiper design which offers more surface contact, and as the plug is pumped down the casing, wiping efficiency is increased. As a top cup on the plug wears, the pressure is transferred to a bottom cup which prolongs the surface engagement maintaining the wiping, resulting in extended wear."}
{"text": "In subsea offshore well operations, a subsea installation may include hydraulically and also electrically controlled equipment, such as valves which are remotely operated from a platform structure located at the surface of the ocean. Hydraulic and electrical control lines extending from the platform to the subsea installation are desired to be releasably connected to the subsea installation so that if the platform must move to another location or be released from its connection to the subsea station, such release may be readily made. In some instances, such a control line has been associated with tension line means integral therewith so that the line, sometimes called an umbilical line, may be locked at its lower end to the subsea installation.\nAn example of such an umbilical line construction is shown in my copending application Ser. No. 771,799 owned by a common assignee wherein the umbilical line extends from the subsea installation to a buoyant structure located about 300 feet below the sea surface and the umbilical line continues above the buoyant structure to a floating platform. The buoyant structure is connected with a marine riser system having a lower riser portion and an upper riser portion so that the upper riser portion and the upper portion of the umbilical line can be disconnected at the buoyant structure and the floating platform moved to another location upon such disengagement. The buoyant structure and the lower riser portion are free-standing to permit reengagement of the upper riser portion at a later time. It is desirable that the buoyant structure be securely tethered to the subsea installation so that in the event the lower riser portion should break, the buoyant structure, which exerts an upward tension on the lower riser portion, will not be suddenly released and rapidly travel to the surface and perhaps strike and cause damage to the floating platform normally thereabove.\nIn prior proposed locking systems for guide line connectors for use in such installations, manually operated locking devices were maintained in locked position by mechanical springs and were unlocked by use of a submarine or by use of go devil means."}
{"text": "The present invention relates to a method for producing an electroluminescent element (may be referred to as an abbreviation of EL) in which a pattern is formed.\nEL elements attract attention as spontaneously emitting flat panel display elements in which positive holes and electrons, which have been poured into from opposed electrodes, are combined in a luminous layer to generate energy, and the energy excites fluorescent material in luminous layer to emit light of a color corresponding to the fluorescent material. Among them, an organic thin film EL display using an organic material as a luminescent material has such a high luminous efficiency that high intensity luminescence can be realized even if applied voltage is a little less than 10 V, is possible to emit light with a simple element structure, and is expected to be applied to low-cost simple expression displays such as advertisements and the like in which specific patterns are emitted and displayed.\nIn the production of a display using such an EL element, patterning on an electrode layer and an organic EL layer is usually put into practice. Methods of EL element patterning include a vapor deposition method of a luminescent material via a shadow mask, a method of separately painting through inkjet, a method of destroying specific luminescent pigment through ultraviolet irradiation, a screen printing method, and others. However, these methods could not provide an EL element being possible to realize all of the high luminous efficiency, the high light takeout efficiency, the simplicity of the production process, and the formation of highly fine patterns.\nThe present invention has been achieved in order to solve the above-mentioned problems. It is a main object of this invention to provide a method for producing an EL element for realizing the high luminous efficiency, the high light takeout efficiency, the simplicity of the production process, and the formation of highly fine patterns."}
{"text": "1. Field of the Invention\nThe present invention relates to an image decrypting apparatus, an image encrypting apparatus, and an image decrypting method that decrypt an encrypted image using a decryption key.\n2. Description of the Related Art\nLeakage of confidential information is a serious problem in the advancement of computerization of the general society and development of a technique that prevents the leakage of the information is desired. For example, as to digital data, a technique has been developed of encrypting data to prevent the content of the data from being viewed when the data is delivered to a third party. The technique is already utilized as a useful means of preventing information leakage.\nOn the other hand, a technique of preventing information leakage of a printed item printed on a paper medium, etc., has not been fully developed yet and any example thereof is not present. It is also actually said that substantially a half of the information leakage occurs from printed items. Therefore, similarly to the case of digital data, development of a technique of preventing information leakage from printed items is an urgent mission.\nSpecific examples for which a measure against the information leakage from the printed items is desired include a bill for purchase of merchandise, a statement for a credit card, etc., a clinical chart of a hospital, a report card of a school, and a name list. The present invention is usable as a technique of preventing information leakage by encryption and decryption of important portions of, for example, each of these examples.\nConventionally, a known example that handles encryption of a printed item can be, for example, a technique disclosed in Japanese Patent Application Laid-Open Publication No. 8-179689. According to this '689 publication: an entire original image is divided into a plurality of blocks; images respectively of the divided blocks are rearranged based on a parameter obtained from a password input (encryption key); the image of each of the blocks designated by the parameter are black/white-inversed and mirror-inversed; and, thereby, the original image is encrypted. When the encrypted image is decrypted, a frame for positioning is added to the image in the outer periphery thereof and another password (decryption key) is input. Thereafter, the original image is recovered by performing the encryption in a reverse manner.\nAnother convention technique can be, for example, a technique of embedding binary data as an image in a printed item, as shown in Japanese Patent Publication No. 2938338. The visualization as an image in this conventional technique is realized by representing the binary data using black-and-white squares each having a designated size and arranging the squares in a matrix. In addition, the printed item is added with a symbol for positioning at a designated position in the matrix such that the position that is visualized as an image can be known in decryption. Using this positioning symbol as the criterion enables shooting of the image using a scanner or a camera and, thereby, decryption of the information embedded.\nIn the above conventional techniques, a process is performed of decrypting an encrypted image output on an analog medium by printing, etc., into the original image using a reverse procedure of the encryption. However, when the contour of the encrypted image is smeared due to smears of ink and optical deformation resulted from printing a digital image and reading the image using an optical apparatus such as a camera or a scanner, the original pixel values can not be recovered even when the reverse process is performed. Therefore, a problem has arisen that the image quality of the decrypted image is degraded.\nAs to the image in the encrypted area, this image is encrypted using a method of dividing the image into a plurality of blocks and rearranging (black/white inversion, mirror inversion, and upside-down inversion) the blocks. However, it can be considered that the image present before the encryption is speculated from the feature of the image of each block of the divided blocks when the blocks are only rearranged. Therefore, a problem has arisen that the security needs to be enhanced."}
{"text": "The present invention comprises a new and distinct variety of poinsettia, botanically known as Euphorbia pulcherrima and hereinafter referred to by the variety name ‘NPCW14233’. ‘NPCW14233’ is a selection from a branch mutation from the poinsettia variety ‘NPCW10158’ (U.S. Plant Pat. No. 22,580).\nIn November/December of 2008 a branch mutation was discovered on the poinsettia variety ‘NPCW10158’ in Nairobi, Kenya. Vegetative cuttings of the branch mutation were taken and propagated. In October of 2009, a single plant was selected in Stuttgart, Germany. In May of 2010, ‘NPCW14233’ was first asexually propagated by apical vegetative cuttings. ‘NPCW14233’ was found to reproduce true to type in successive generations of asexual propagation via apical vegetative cuttings in Stuttgart, Germany."}
{"text": "When designing a handset for use in mobile telephony, handset designers generally start with an ear model. Additionally, handset designers typically pay close attention to the requirements of the mobile telephone network operators (carriers) that will market and sell the handsets. Currently, carriers require that handsets be designed based on either a “high-leak” ear model or a “low-leak” ear model. In the high-leak ear model, an attempt at forming a seal between a mobile telephone handset and a user's ear is expected to form a consistent seal that has a significant acoustic leak. In the low-leak ear model, an attempt at forming a seal between a mobile telephone handset and a user's ear is expected to form a consistent seal that has insignificant acoustic leak.\nA user of that handset may change the position of their ear relative to the handset, and subsequently the acoustic leak, throughout a phone call or based on personal holding position. Also, each user has a different acoustic load due to variations in ear size and shape; accordingly, the user may prefer a lower leak design or a higher leak design.\nTuning a handset often involves creating an optimum acoustic equalizer (EQ) setting. An acoustic EQ setting determines which frequencies of received audio are amplified, and to what extent, before the received audio is sent to a handset speaker, and, potentially, which frequencies are suppressed, and to what extent. The goal is to meet a specific “frequency-versus-level mask”. Such masks are determined by standards bodies, such as The 3rd Generation Partnership Project (3GPP, see www.3gpp.org).\nHowever, during operation, users are often unable to form a consistent seal between the ear and the handset, due to movement and changes in position during the voice call. The amount of acoustic leak from the seal is dependent on the manner in which the handset is held by the user and the size and shape of the user's ear. Under those circumstances wherein the user forms a low-leak seal and the handset has been designed with a single acoustic EQ setting that assumes a high-leak seal, the user may complain about poor audio quality. For example, the user may complain that received audio sounds “boomy” or “muddy”. Given the variety of user ear shapes and sizes, along with the variety of manners in which users hold the handset, it may be quite difficult to establish a single acoustic EQ setting that will meet the requirements of all users and certification bodies."}
{"text": "1. Field of the Invention\nThe present invention relates generally to carbohydrate biochemistry. More specifically, the invention relates to starch biosynthesis and the enzyme(s) involved.\n2. Description of the Related Art\nStarch, the most significant carbohydrate reserve in plant storage tissues, comprises the glucose homopolymers amylose and amylopectin. Amylose consists of predominantly linear chains of xcex1-(1xe2x86x924)-linked glucose residues, whereas amylopectin is a highly branched glucan with a specific xe2x80x9cclusteredxe2x80x9d distribution of xcex1-(1xe2x86x926) glycosidic bonds (i.e. branch linkages) connecting linear chains (French, 1984; Manners, 1989).\nDespite the relatively simple chemical structure of amylopectin, very little is known about the enzymatic processes responsible for formation of the highly specific and complex branching patterns in this polysaccharide. Biosynthesis of amylose and amylopectin involves activities of four groups of enzymes, each of which comprises multiple isozymes. These enzymes are ADPG pyrophosphorylases (AGPase), starch synthases (SS), starch branching enzymes (SBE) and starch debranching enzymes (SDBE) (Preiss, 1991; Hannah et al., 1993; Martin and Smith, 1995; Nelson and Pan, 1995; Ball et al., 1996; Preiss and Sivak, 1996; Smith et al., 1996). These enzymatic steps can account for all chemical linkages in starch, however, the specific roles of individual isozymes in formation of specific branching patterns in amylopectin and determination of starch granule structure and properties remain unknown.\nAnalysis of maize mutants with abnormal endosperm phenotypes has contributed greatly to the understanding of starch synthesis (Shannon and Garwood, 1984; Nelson and Pan, 1995) and facilitated the identification of many genes coding for starch biosynthetic enzymes. Cloned genes whose products are thought to be involved directly in starch biosynthesis are waxy (wx), coding for the granule-bound starch synthase GBSSI (Shure et al., 1983; Klosgen et al., 1986), amylose extender (ae), coding for SBEIIb (Fisher et al., 1993; Stinard et al., 1993), shrunken2 (sh2) and brittle2 (bt2), coding for the large and small subunits of AGPase, respectively (Bae et al., 1990; Bhave et al., 1990), and sugary1 (su1), coding for the SDBE SU1 (James et al., 1995). The transposon-tagging strategy was used to determine that the abnormal endosperm phenotype of wx-, ae-, or su1-mutants results from primary defects in GBSSI, SBEIIb, or SU1, respectively, and this approach remains the most effective way to identify genes such as dull1 (du1), in which the primary defect can not be associated with a particular enzyme deficiency.\nThe du1-mutations define a gene with a very important function in starch synthesis, as indicated by extensive structural analyses of starch from du1-mutant endosperms, and by the effects of these mutations when combined with other genetic deficiencies in starch biosynthetic enzymes (Shannon and Garwood, 1984; Nelson and Pan, 1995). The reference mutation du1-Ref was first identified as a recessive modifier of su1-Ref and su1-amylaceous (su1-am) (Mangelsdorf, 1947). Mutations of du1, when homozygous in otherwise non-mutant backgrounds, result in mature kernels with a tarnished, glassy, and somewhat dull appearance referred to as the xe2x80x9cdull phenotypexe2x80x9d. Expression of this phenotype, however, depends on the particular genetic background (Mangelsdorf, 1947; Davis et al., 1955). Total carbohydrate and starch content in mature du1- mutant kernels is slightly lower than normal (Creech, 1965; Creech and McArdle, 1966). The apparent amylose content in starch from du1-mutants is slightly or greatly elevated compared to normal depending on the genetic background (Shannon and Garwood, 1984), although the properties of polysaccharides in the apparent amylose fraction are essentially not altered (Dvonch et al., 1951). Approximately 15% of the starch in du1-mutant endosperms is in a form known as xe2x80x9cintermediate materialxe2x80x9d, which is distinguished from amylose and amylopectin by the properties of its starch-iodine complex (Wang et al., 1993b). Analysis of combined amylopectin/intermediate material fractions indicated that starch from du1-mutants has the highest degree of branching among a wide variety of normal and mutant kernels analyzed (Inouchi et al., 1987; Wang et al., 1993a; Wang et al., 1993b). Starch granules from du1-mutants seem to have normal structural and physical properties, although some abnormally shaped granules are found in the mutant endosperm (Shannon and Garwood, 1984).\nDespite these subtle effects exerted by the single mutation, du1-alleles when combined with other mutations affecting starch synthesis result in a broad range of more severe alterations (Shannon and Garwood, 1984; Nelson and Pan, 1995). Mutations of du1 have been examined in combination with wx-, ae-, su1-, and sugary2 (su2-) mutations, and in all instances the double mutant kernels contained more soluble sugars and less total starch than when any of the mutations was present alone. In many instances the double mutants also produce polysaccharide forms that are distinct from the starch found in any single mutant kernels. These pleiotropic effects indicate the product of Du1 affects many aspects of starch biosynthesis in maize endosperm, however, without knowing the identity of this protein it is difficult to assess its specific functions.\nConsistent with the pleiotropic genetic effects, du1-mutations cause reduced activity in endosperm of two seemingly unrelated starch biosynthetic enzymes, the starch synthase SSII and the branching enzyme SBEIIa (Boyer and Preiss, 1981). SSII is one of two enzymatically distinct starch synthase activities identified in the soluble fraction of maize endosperm; in vitro activity of SSII requires an exogenous glucan primer, and its molecular weight was determined in different studies as either 95 kD or 180 kD (Boyer and Preiss, 1981; Mu et al., 1994). Similarly, SBEIIa is one of the three known SBE isozymes in endosperm cells (Boyer and Preiss, 1978b; Fisher et al., 1993; Fisher et al., 1995; Gao et al., 1997). Several possibilities exist to explain the dual biochemical effects of du1-mutations. Du1 may code for a protein regulating the expression or activity of both SSII and SBEIIa. Alternatively, Du1 may code for either of these two enzymes, and the deficiency in one enzyme might also affect the second enzyme because of a direct or substrate-mediated physical interaction.\nDU1 codes for a starch synthase, as indicated by the extensive-similarity of its deduced amino acid sequence to potato SSIII, and by the substantial similarity between the C-terminal residues of DU1 and a large group of phylogenetically diverse starch- and glycogen synthases. Particularly striking are two regions that together comprise more than half of the deduced DU1 sequence of 1,674 residues, which share very high similarity of 51% and 73%, respectively, with the corresponding regions of the potato SSIII sequence. Within a stretch of 450 amino acids at the C-terminus of DU1 nearly 30% of the best aligned residues are identical in comparisons to a wide variety of starch- and glycogen synthases, suggesting the location of a domain within DU1 that provides xcex1-1,4-glycosyltransferase activity.\nThe starch synthase coded for by Du1 is the soluble isozyme identified biochemically as SSII (Ozbun et al., 1971; Boyer and Preiss, 1981). The deduced molecular weight of DU1 including a potential transit peptide, 188 kD, matches closely with that of 180 kD reported for mature SSII lacking a transit peptide (Mu et al., 1994). The size difference of approximate 8 kD may be due to the transit peptide present in the deduced DU1 sequence. The tissue specific expression pattern of the Du1 mRNA also matches the expression pattern of SSII. Du1 transcripts were undetectable in leaves either by RNA gel blot or RT-PCR analyses, corresponding with that fact no detectable SSII activity was present in leaf extracts (Dang and Boyer, 1988). Moreover, the activity of SSII, along with that of SBEIIa, was greatly reduced in du1-mutant endosperm (Boyer and Preiss, 1981). Therefore, it appears that the maize du1 locus codes for the soluble starch synthase SSII, the counterpart of potato SSIII.\nThis characterization of DU1 implies that the phenotypic effects of du1-mutations, including changes in starch structure, deficiencies of two starch biosynthetic enzymes, and genetic interactions with ae-, su1-, su2-, and wx- mutations, all result either directly or indirectly from alteration of SSII. The reduction of SBEIIa activity in du1-mutant endosperm could result from the SSII deficiency owing to physical interaction between the two enzymes. A direct physical association of SSII and SBEIIa is implied by the observation that peak activities of both SSII and SBEIIa always coincide in the same DEAE-cellulose column fractions (Boyer and Preiss, 1978a; Boyer and Preiss, 1981; Dang and Boyer, 1988). Thus, SSII and SBEIIa may function together in vivo in the form of single multi-enzyme complex. Loss of the intact enzyme complex owing to reduction of SSII in du1-mutant endosperm may result in abnormally rapid proteolytic turnover of SBEIIa, or prevent accumulation of the enzyme by some other mechanism. Alternatively, expression of the Sbe2a gene in du1-mutant endosperm may be inhibited as a more indirect consequence of the deficiency in SSII, for example through reduction of a transcriptional inducer or elevation of a repressor. Although the du1-Ref mutation does indirectly affect expression of other starch biosynthetic genes (Giroux et al., 1994), it actually caused increased gene expression rather than the reduction observed for SBEIIa. Furthermore, considering that large glucose polymers are expected to be the substrate and product of DU1, down-regulation of Sbe2b expression by a transcriptional mechanism seems unlikely. Thus, the former hypothesis may explain the deficiency of SBEIIa in du1-mutant endosperm.\nThe broad impact of the combination of du1-mutations with various su1-alleles on kernel phenotype and starch synthesis (Cameron, 1947; Shannon and Garwood, 1984) could be explained by the SU1 SDBE also interacting closely with SSII in vivo, perhaps in the same enzyme complex with SBEIIa. This proposed association of SBEIIa and SU1 in a multi-enzyme complex is consistent with the proposed simultaneous branching and debranching actions during amylopectin synthesis by SBE and SDBE (James et al., 1995; Nelson and Pan, 1995; Ball et al., 1996).\nThus, the prior art is deficient in understanding the complex association of enzymes involved in starch synthesis and in cloning genes corresponding to these enzymes. The present invention fulfills this long-standing need and desire in the art.\nTo illustrate the role of the du1locus in starch biosynthesis, a transposon-tagging strategy was used to isolate the gene and describe its polypeptide product. The present invention reports tagging of the du1 locus with a Mutator (Mu) transposon, cloning and characterization of a portion of the gene, and the sequence of a near full-length cDNA (SEQ ID No. 1). The amino acid sequence deduced from this cDNA indicates Du1 codes for a 186 kD polypeptide extremely similar to SSIII, a starch synthase from potato tubers (Abel et al., 1996; Marshall et al., 1996). The expression pattern of Du1 also was characterized. Taken together these characterizations indicate that Du1 most likely codes for SSII of maize endosperm. In addition, the product of Du1 contains unique sequence features in its amino terminus that may mediate direct interactions with other starch biosynthetic enzymes.\nOne object of the present invention is to provide an enzyme with which to regulate the production of starch, and with which to produce altered or novel forms of starch.\nIn an embodiment of the present invention, there is provided a cDNA corresponding to the dull1 gene of maize.\nIn yet another embodiment of the present invention, there is provided an expression vector containing the sequence of dull1 with which to produce the starch synthase enzyme in transgenic plants or other prokaryotic or eukaryotic organism.\nIn yet another embodiment of the present invention, there is provided (1) cDNA having the nucleotide sequence comprising nt 120 to nt 1221 of SEQ ID No. 1, said sequence encoding the first 368 amino acids of DU1; (2) cDNA having the nucleotide sequence comprising nt 655 to nt 1221 of SEQ ID No. 1, said sequence encoding amino acids 180 to 368 of DU1; (3) cDNA having the nucleotide sequence comprising nt 565 to nt 816 of SEQ ID No. 1, said sequence encoding amino acids 150 to 233 of DU1; (4) cDNA having the nucleotide sequence comprising nt 1369 to nt 1944 of SEQ ID No. 1, said sequence encoding amino acids 418 to 609 of DU1; (5) cDNA having the nucleotide sequence comprising nt 1 to nt 1437 of SEQ ID No. 1, said sequence encoding amino acids 1 to 440 of DU1; (6) cDNA having the nucleotide sequence comprising nt 1438 to nt 2424 of SEQ ID No. 1, said sequence encoding amino acids 441 to 769 of DU1; (7) cDNA having the nucleotide sequence comprising nt 2425 to nt 3791 of SEQ ID No. 1, said sequence encoding amino acids 769 to 1225 of DU1\nOther and further aspects, features, and advantages of the present invention will be apparent from the following description of the presently preferred embodiments of the invention. These embodiments are given for the purpose of disclosure."}
{"text": "Backlight is an important constituent part of a liquid crystal module. A main operation module of a backlight is a Light Emitting Diode (LED) driver and an LED light. Therefore, an operation voltage and operation current of the LED light are major factors which influence power consumption of the backlight.\nA calculation formula of the power consumption of the backlight is: power consumption=voltage×current. The LED light is a physical element with PN junction characteristics, and as operation time elapses, current flowing through the LED light may increase with temperature increases at a constant voltage, and thus the power consumption also increases therewith."}
{"text": "The present invention relates to a phase change memory.\nPhase change memories use a class of materials that have the property of switching between two phases having distinct electrical characteristics, associated to two different crystallographic structures of the material, and precisely an amorphous, disordered phase and a crystalline or polycrystalline, ordered phase. The two phases are hence associated to resistivities of considerably different values.\nCurrently, the alloys of elements of group VI of the periodic table, such as Te or Se, referred to as chalcogenides or chalcogenic materials, can be used advantageously in phase change memory cells. The currently most promising chalcogenide is formed from an alloy of Ge, Sb and Te (Ge2Sb2Te5).\nIn the chalcogenides, the resistivity may vary by two or more orders of magnitude when the material passes from the amorphous (more resistive) phase to the crystalline (more conductive) phase, and vice versa.\nPhase change can be obtained by locally increasing the temperature. Below 150° C., both the phases are relatively stable, with an amorphous bit tending to gradually crystallize over a period of years if held at 150° C. Starting from an amorphous state, and raising the temperature above 200° C., where there is a rapid nucleation of the crystallites and, if the material is kept at the crystallization temperature for a sufficiently long time, it undergoes a phase change and becomes crystalline. Preferably, the material is cooled off slowly to better assure the crystalline state. To bring the chalcogenide back to the amorphous state it is necessary to raise the temperature above the melting temperature (approximately 600° C.) and then rapidly cool off the chalcogenide."}
{"text": "1. Field of the Invention\nThe present invention relates to a device which separates various objects from image data and which performs image processing on the objects thus extracted, and more specifically relates to a device which determines whether to perform the vector transform on an object in accordance with a characteristic of the object, i.e. an object attribute and/or the type of job set for the object.\n2. Description of the Related Art\nThe spreading of digital cameras and image scanners today allows photo images to be digitalized easier than before. What has been made possible, for example, is: taking digitized photo images into a personal computer (hereafter, simply referred to as a PC); storing the images as image files; and processing and editing the image files by application programs. In addition, the image files of the photo images thus processed and edited can be laid out in an electronic document, such as a text document and a drawing. In this way, an electronic document that has photo images incorporated therein can be created.\nIn addition, a digital multifunctional peripheral that integrates functions of a scanner, a copier, a printer, a fax, and the like can be connected to a PC through a communication network, such as a LAN and the Internet. Such networking allows the PC to send, to the digital multifunctional peripheral, such electronic documents as a text document or a drawing incorporating a photo image taken by a digital camera. Then, the digital multifunctional peripheral can print out the received electronic document.\nThere exists a well-known image processing system capable of compressing image data to a higher degree while maintaining quality of an image, by separating the image data into a plurality of objects, and then by performing image processing and encoding processing on each of the objects. In addition, Japanese Patent Laid-Open No. H05-211607 discloses a method of reducing image degradation. The method includes: extracting a character object from image data; detecting a contour of the character object; and switching between smoothing the detected contour and not smoothing the detected contour. Here, whether to perform smoothing the detected contour is determined in accordance with a magnification specified by a user.\nIn the above-described conventional technique, after a character object is extracted from image data, transform processing from the image data into vector data (a vector image), such as a contour-extraction processing is always performed without consideration of a characteristic of the extracted character object. Note that, the contour-extraction processing, in general, takes a large amount of time. Accordingly, transforming image data into vector data without considering the characteristic of the object is likely to be a factor to lower the performance of the device."}
{"text": "1. Technical Field\nThe present disclosure relates to a lens barrel provided with an optical system.\n2. Related Art\nJP 2012-53411 A discloses a lens barrel which enables zooming."}
{"text": "Enteropeptidase is a serine protease that converts trypsinogen secreted from the pancreas after meal to trypsin. Trypsin in a state activated by enteropeptidase then activates protease precursors such as chymotrypsinogen, procarboxypeptidase, and proelastase. These activated proteases decompose dietary proteins into amino acid units. The resulting amino acids are absorbed into the small intestine. Thus, enteropeptidase inhibitors are capable of suppressing the degradation or absorption of proteins and is useful as a drug for treating obesity.\nExamples of heterocyclic compounds include the following:\n(1) A compound which has a trypsin inhibitory effect and is useful in the treatment or prevention of a renal disease or a disease involving trypsin, the compound being represented by the following formula:\nwherein\nring A represents\n\nR1 represents halogen, lower alkyl, or OH;\nR2 represents H, optionally substituted cycloalkyl, optionally substituted aryl, an optionally substituted aromatic heterocyclic ring, an optionally substituted nonaromatic heterocyclic ring, —C(O)-lower alkylene-optionally substituted aryl, or lower alkyl optionally substituted by halogen or the like;\nL1 represents —Y1-lower alkylene-Y2—, or —C(O)—N(R6)—;\nY1 represents a bond or —C(O)—;\nY2 represents a bond, —N(R6)—, or —C(O)—N(R6)—;\nL2 represents -(lower alkylene optionally substituted by CO2H or the like)-, —Y3-cyclohexanediyl-Y4—, or —Y3-phenylene-Y4—, and L2 optionally forms optionally substituted cyclic amino together with R2;\nY3 represents a bond or lower alkylene;\nY4 represents a bond, lower alkylene, or —C(O)—;\nR3 represents H, lower alkyl optionally substituted by halogen, halogen, OH, —O-lower alkyl, cycloalkyl, aryl, or the like;\nR4 represents lower alkyl optionally substituted by halogen, halogen, OH, —O-lower alkyl, cycloalkyl, aryl, or the like;\nR5 and R6 each represent H or lower alkyl;\nX1, X2, and X3 each represent CH or N (except that at least one of X1, X2, and X3 is N);\nm represents an integer of 0 to 4;\np represents an integer of 0 to 3; and\nq represents an integer of 0 to 4\n(Patent Literature 1).\n(2) A compound which has a serine protease inhibitory effect and is useful in the treatment or prevention of obesity, hyperlipidemia, diabetes mellitus, diabetic complications, or metabolic syndrome, the compound being represented by the following formula:\nwherein\nR1, R2, R3, and R4 each represent H or the like;\nHetAr represents an optionally substituted heteroaromatic ring;\nX represents optionally substituted lower alkylene or the like;\nY represents carbonyl or the like;\nA represents\nor the like; and\nR6 and R7 each represent H, optionally substituted lower alkyl, or the like\n(Patent Literature 2).\n(3) A compound which has a serine protease inhibitory effect and is useful in the treatment or prevention of obesity, hyperlipidemia, diabetes mellitus, diabetic complications, or metabolic syndrome, the compound being represented by the following formula:\nwherein\nD represents a benzene ring, a naphthalene ring, or a pyridine ring;\nHet represents a heterocyclic ring;\nR1 represents H or the like;\nR2 represents nitro, lower alkyl, or the like;\nX represents optionally substituted lower alkylene;\nZ represents —N(R3)-(wherein R3 represents H, optionally substituted lower alkyl, optionally substituted lower cycloalkyl, or the like);\nY represents a single bond or —(CH2)p-C(R4a) (R4b)-(CH2)q- (wherein R4a and R4b each represent H, lower alkyl, or aralkyl, and p and q each represents an integer of 0 to 5); and\nA represents —CO2R6 (wherein R6 represents H or lower alkyl) or\nwherein Q represents optionally substituted lower alkylene, and R7 represents H or lower alkyl(Patent Literature 3).\n(4) A compound which has an enteropeptidase inhibitory effect and is useful in the treatment or prevention of a disease related to obesity or abnormal fat metabolism, the compound being represented by the following formula:\nwherein\nB represents boron;\nW represents a nitrogen-containing functional group (\n\nX represents a linker (CX1X2)p;\nY and Z each represent OH, OR (R represents alkyl), a homocyclic ring, a heterocyclic ring, or the like;\nR2 represents aminoacyl, acyl, or the like; and\nR2 represents H, alkyl, or OR (R represents H or alkyl)\n(Patent Literature 4).\n(5) A compound which has a serine protease inhibitory effect and is useful in the treatment or prevention of obesity, diabetes mellitus, or the like, the compound being represented by the following formula:\nwherein\nR1 and R2 each represent alkyl or the like; and\nX represents —OR3, —NR4R5, or the like\n(Patent Literature 5)."}
{"text": "It is already a known practice to place turbulators in the fluid circulation tubes of heat exchangers in order to increase heat exchanger efficiency by forming turbulence in the flow of fluid along the tubes.\nUp to now, turbulators have often been in the form of highly elongate, very thin, rectangular strips of ductile material such as aluminum having baffles punched out from the strip and projecting from both faces thereof. However, it to be observed that such turbulators are deformable and difficult to insert into the tubes, in particular when the tubes are of oblong or oval cross section. This is because the turbulator strips are so thin and because of the baffles projecting from both sides thereof.\nPreferred embodiments of the present invention mitigate these drawbacks."}
{"text": "Cut flowers have become increasingly popular and are sold at florists, in grocery stores or the like, and often by street vendors. One of the disadvantages of buying cut flowers at the less expensive outlets, i.e., street vendors and/or supermarkets, is that they are often delivered to you, and you hence deliver them, in unattractive plastic bags.\nA viable alternative to the plastic bag is a, heretofore unavailable, attractive vase structure which is inexpensive and may be stored in a flat condition, but quickly and easily assembled into a water-and-stem-holding container.\nPrior art known to the inventor, in addition to the plastic bag method described hereinabove, includes U.S. Pat. No. 404,585 granted to Wright, June 4, 1889, which discloses a pot made of paper having its bottom folded inwardly such that once the pot, containing dirt and rootball are positioned in the appropriate place beneath the soil, the pot is pulled upwardly providing a protective shield for the seedling.\nU.S. Pat. No. 2,150,453 granted to Mulford, et al., Mar. 14, 1939, discloses several packages which may be constructed from precut and prefolded blanks to form a flower pot. A bulb, purchased in the pot, may be grown in the pot by the purchaser such that the bulb may never need to be transplanted.\nU.S. Pat. No. 2,340,373 granted to Gardner, Feb. 1, 1944, discloses a plant pack that includes a single-piece exterior shipping container and interior means to secure the plant within the container preventing movement during shipping.\nU.S. Pat. No. 3,373,917 granted Mar. 19, 1968 to Cox, discloses a cup including a foldable exterior casing and a fluid-impervious liner.\nU.S. Pat. No. 4,209,092 granted June 24, 1980 to Jones, discloses a plant container fabricated of a single sheet of material having predetermined fold lines and including a sealing strip for securing the container in its usable three-dimensional configuration.\nU.S. Pat. No. 4,347,686 granted Sept. 7, 1982 to Wood, discloses a container for plant life which is folded such that it forms a cone or pyramid for facilitating aerial planting. The device may include a solid nose piece and fins for guidance."}
{"text": "Current high-volume, high-availability, content distribution networks use clusters of Web servers to achieve scalability and reliability. To serve a large and diverse client population, content can be replicated across servers, or partitioned with a dedicated server for particular content or clients. In such environments a front-end dispatcher (usually called a “web switch” or a “reverse proxy” based on its functional usage) directs incoming client requests to one of the server machines. The request-routing decision can be based on a number of criteria, including requested content, server load, client request, or client identity.\nDispatchers are typically required to perform several functions related to the routing decision. These include:                distribute incoming requests as to balance the load across servers;        examine client requests to determine which server is appropriate to handle the request (i.e., content-based routing);        identify the client to maintain affinity with a particular server for e-business applications or to provide service differentiation.        \nDispatchers may be broadly categorized into two types: layer-4 dispatchers which base decisions on TCP and IP headers alone, and layer-7 dispatchers which use application layer information (e.g., HTTP headers) to direct clients. The use of layer-4 or layer-7 dispatchers depends on the request-routing goal. Load-balancing across replicated content servers, for example, typically does not require knowledge of the client request or identity, and thus is well-suited to a layer-4 approach. Content-based routing requires the dispatcher to terminate the incoming connection, examine the higher-layer headers, and either create a new connection with the appropriate server (using connection splicing), or transfer the connection to the appropriate server (using connection handoff). Layer-7 dispatchers, while more sophisticated, suffer from limitations on scalability and performance since they must perform connection termination and management for a large number of clients.\nHigh-speed switching hardware is also commonplace in core ISP networks with a growing migration to Multiprotocol Label Switching (MPLS). MPLS provides a circuit-switching service in a hop-by-hop routed network and is presently used for flexible routing, traffic engineering and Virtual Private Networks (VPNs). It achieves this by grouping related packets by assigning them a common, fixed-size label. Packets sharing a label belong to the same forwarding equivalence class (FEC) and can be routed and treated the same way in the network. Standard usage of MPLS involves establishment of arbitrary label-switched paths (LSPs) for forwarding particular classes of traffic. LSPs may also be nested by stacking MPLS labels where an outer label might be used to assign traffic to a common network-wide path, while an inner label could be used to demultiplex traffic among classes of traffic on that path. An MPLS-enabled network includes of label-switched routers(LSRs) that implement the MPLS protocols.\nMPLS labels do not have built-in semantics, i.e. labels are used only to map a packet from an LSR input port to an LSR output port. It would be advantageous to make use of this flexibility by mapping application-layer information onto labels to enable high-performance Web switching, rather than using labels to express routing and forwarding policies (as is customary)."}
{"text": "1. Field of the Invention\nThe present invention relates to metal oxide based memory devices and methods for manufacturing such devices.\n2. Description of Related Art\nResistive random access memory, RRAM, is a type of nonvolatile memory that provides the benefits of small cell size, scalability, ultrafast operation, low-power operation, high endurance, good retention, large On/Off ratio, and CMOS compatibility. One type of RRAM includes metal oxide layers than can be caused to change resistance between two or more stable resistance ranges by application of electrical pulses at levels suitable for implementation in integrated circuits.\nIn prior methods of forming a transition metal oxide (TMO) memory layer, photoresist is used as an etching mask to protect the portions of TMO in a designated region that will be the memory layer, while other portions of TMO that are not part of the memory layer are etched away. After etching, the photoresist is stripped and this stripping may cause damage to the TMO memory layer. Prior methods of avoiding this damage caused by stripping have included using a hard mask such as SiO2 or Si3N4 prior to oxidization to define surfaces to be oxidized to form the TMO memory layers. By defining which regions to oxidize, the step of etching away unwanted oxidized regions is eliminated. However, this process of using a hard mask adds two extra processes to the previously described method involving damage caused by the photoresist strip. These additional steps in the process add time and cost to production, which are undesirable.\nIt is therefore desirable to provide a memory cell and method of manufacture that eliminates the possibility of damage to the TMO memory layer caused by a photoresist strip while adding a minimal amount of manufacturing steps in order to provide a cost-effective method of manufacture."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of forming an oxide film and further a method of manufacturing an electronic device utilizing an oxide film.\n2. Description of the Prior Art\nIn recent years, researchers are attracted by thin film transistors utilizing non-single crystalline semiconductor thin films.\nConventionally, such a non-single crystalline semiconductor thin film is formed on an insulating substrate by chemical vapor deposition, so that a temperature during the film formation is as low as 450° C. or less. Therefore, soda-lime glass, boro-silicate glass, and the like can be used as the substrate.\nThe thin film transistor recently attracting researchers is a field effect transistor (simply referred to as FET) having the same function as that of MOS FET. The size of the thin film transistor is limited only by the size of the apparatus to be used for formation of a semiconductor thin film constituting the transistor, so that it is easy to form transistors on large-sized substrates. Such large-sized thin film transistors are promising. For example, the large-sized thin film transistors can be used as switching elements of liquid crystal displays having a lot of pixels in the form of matrix or switching elements of one dimensional or two dimensional image sensors or the like.\nIt is possible to implement a conventional fine processing to the semiconductor thin films. Hence, the thin film transistor can be formed by means of a conventional fine processing, for example photolithography technique. And it is also possible to make the thin film transistor integrated as a function element of a part of monolithic IC.\nReferring to FIG. 2, a typical structure of a conventional thin film transistor is schematically illustrated.\nSource and drain electrodes 24 and 25 are provided on an insulating substrate 20 made of glass and source and drain regions 22 and 23 are provided on the source and drain electrodes 24 and 25 respectively and a non-single crystalline semiconductor thin film 21 is provided on the substrate 20 and a gate insulating film 26 is provided on the semiconductor thin film 21 and a gate electrode 27 is provided on the gate insulating film 26.\nIn the thin film transistor, electric current flowing between the source region 22 and the drain region 23 is controlled by a voltage applied to the gate electrode 27.\nA gate oxide film constituting such a thin film transistor was conventionally formed by exposing a semiconductor material to thermal oxidation or by thermal CVD under a reduced or atmospheric pressure, or the like.\nElectric characteristics of the thin film transistor largely depend on the quality of a channel region of the semiconductor film and the quality of the gate insulating film. For this reason, a gate insulating film of particularly good quality has eagerly been required.\nIn the case of the formation of the gate oxide film by exposing a semiconductor material to thermal oxidation or by thermal CVD under a reduced or atmospheric pressure, the temperature during the formation of the gate insulating film should be as high as approximately 600° C. in order to obtain a thin film transistor of good electric characteristics. So that, a heat-resistant substrate material such as quartz glass had to be utilized though it is expensive.\nWith respect to a method for forming a gate insulating film at a low temperature, a plasma CVD and a sputtering method utilizing an argon gas for sputtering are well-known. This sputtering method is implemented in an atmosphere comprising a large amount of argon, specifically an atmosphere comprising 100 to 80 volume % Ar atoms and 0 to 10 volume % oxygen. This is because probability of an atom or a cluster of atoms being dislodged from a target by collision of one inert gas atom, for example one Ar atom, is high (in other words, sputtering yield of Ar gas is high). However, in both the plasma CVD and the sputtering method using a large amount of argon, the gate insulating film involves numbers of elements (e.g. inert gas elements such as Ar) which was involved in a target or existed in a chamber during the CVD or the sputtering, resulting in generation of fixed electric charges in the gate insulating film. Further, ions of the elements bombard a surface of an activated layer in a thin film transistor and thereby give a damage thereto. Hereupon, a mixed layer of the activated layer and the gate insulating film is formed in the vicinity of an interface between the activated layer and the gate insulating film. As a result, interfacial level is formed at the interface and a thin film transistor of fine characteristics cannot be obtained by any of those methods.\nIt has been attempted to form a gate insulating film by a photo CVD method, and an interfacial level density of the gate insulating film was about 2×1010 eV−1 cm−2, almost the same as that of a thermal oxidation film. However, the photo CVD method required a long period of time, in other words, the film formation speed was extremely slow, so that the photo CVD method was not suitable for an industrial application.\nReferring now to FIG. 7, a network of silicon oxide formed by sputtering in an atmosphere comprising a large amount of argon is illustrated. Symbols 0 in the drawing indicate oxygen or silicon and symbols X indicate dangling bonds of silicon. A silicon oxide film including a gate insulating film is not dense when quantity of fixed electric charges is large in the silicon oxide film. The larger the number of dangling bonds of silicon in the silicon oxide film is, the larger the quantity of fixed electric charges is. And the larger the number of Ar+ in the silicon oxide film is, the larger the quantity of fixed electric charges is. Ar+ and Ar tend to stay inside the silicon oxide network as illustrated in FIG. 7 (Ar+ and Ar do not tend to be substituted for Si or 0 in the network). In fact, numbers of dangling bonds of silicon tend to be generated in the silicon oxide film when the silicon oxide film is formed by sputtering in an atmosphere comprising a large amount of argon. This is partly because internal stress is generated in the silicon oxide film by Ar or Ar+ present inside the silicon oxide network and partly because defects are formed in the silicon oxide film by bombardment of argon with the silicon oxide film during sputtering."}
{"text": "In base stations for personal communications systems, for example used for the Global System for Mobile communications (GSM), Enhanced Data Rates for GSM Evolution (EDGE), Wideband Code Division Multiple Access (W-CDMA), radio frequency (RF) power amplifiers are among the key components. For these power amplifiers, RF power laterally diffused metal oxide semiconductor (LDMOS) transistors are a commonly used technology. They exhibit excellent high power capabilities, gain and linearity. These metal oxide semiconductor (MOS) transistors are not only used in base stations but also in radar and broadcast applications. If a RF mismatch occurs at the output of the LDMOS, power will be reflected back at the LDMOS and the LDMOS must be able to handle this. Thus ruggedness, i.e. the capacity to sink excess energy without damage, is an important issue for RF power LDMOS transistors. For low-frequency applications, ruggedness can be improved by adding external diodes with appropriate breakdown voltages. The diodes will sink electrical energy and thus prevent damage to the LDMOS. For RF LDMOS transistors with high frequency applications, external diodes would degrade the RF performance and the ruggedness improvement would be insufficient. It is also conceivable to sink electrical energy through the drain to substrate diode, if the drain-substrate diode is able to sink enough energy before the voltage reaches a level at which the parasitic bipolar transistor is turned on and the LDMOS is damaged. However, the breakdown often starts locally in a specific section of the drain, which determines a low breakdown voltage but doesn't give a high breakdown current.\nIt is therefore an object of the invention to improve a laterally diffused metal oxide semiconductor circuit for a radio frequency power amplifier in a way that allows the adjusting of a breakdown voltage at which breakdown will occur with a higher breakdown current."}
{"text": "The present invention relates to coating processes and compositions, and more particularly to processes and coating compositions for suppressing combustion of metals, and to the products of such processes.\nThe relatively high strength, low density and high melting point of titanium has led to increasing use of titanium and titanium alloys in aircraft engine components such as fan and compressor vanes and blades, and in other applications. However, titanium can be made to ignite and react with air in an exothermic self-sustaining rapid oxidation process referred to as combustion. Titanium has a high heat of combustion, low thermal conductivity and a spontaneous ignition temperature below its melting point. These properties favor ignition and subsequent propagation of titanium or titanium alloy combustion rather than melting, thus producing additional rapid local temperature increases and rapid propagation of the resultant combustion once ignition occurs. During operation of aircraft jet engines, titanium blade ignition is often caused by rub. An abradable seal of rubber, felt metal or nickel-graphite, is present on the engine shroud adjacent to the compressor blades to prevent gas backflow and consequent loss of pressure. Without such a seal, efficiency of the compressor would be severely impaired. Accordingly, seal-blade separation is necessarily small. During engine operation, the blades heat up, expand with heat, and even creep, so that some blades thereby rub against the seal. Severe rubbing and abrasion may occur when an engine part fails, putting the engine out of balance, or when a failed blade imbeds itself in the seal. This abrasion between seal and blade(s) can cause ignition of the titanium. The high-velocity air stream in jet engines aids further to the combustion of any titanium or titanium alloy component that ignites and causes burning particles and molten metal to be sloughed off therefrom. These particles can be entrained in the engine air stream and impinge on downstream engine components to ignite them and thereby spread combustion in the engine. Some single-layer metallic and organic coatings have been used for reducing combustion of titanium substrates. However, an organic coating, particularly its binder, will generally vaporize or burn off, before its protection is needed, in a 600.degree. F.-900.degree. F. working environment such as found in operating jet turbine aircraft engines. Also, many metals, particularly those capable of forming an intermetallic with the substrate metal, can if coated on the substrate diffuse into and significantly degrade engineering or mechanical properties, such as fatigue endurance, of the substrate, particularly in a relatively hot gas turbine engine environment."}
{"text": "Electronic navigation systems provide for the identification and tracking of the physical location of a user and are, in many cases, also capable of directing users to a destination of their choice. These systems enable the presentation of a variety of way finding information such as the current physical location of the system user, intermediate waypoints and destination. Navigation systems can use both audio and/or visual methods to guide users to a destination. In addition to stored data relating to the physical environment, these systems may also allow users to locate businesses and points of interest by category, keywords, or through other querying capabilities. The inclusion of a particular business or point of interest in an electronic navigation system may or may not require a business or individual to pay a fee to provide for this inclusion. In addition to passive listings, navigation systems also have the ability to present advertisements to users, potentially allowing them to initiate a trip based on the advertisement.\nIn general, and unlike traditional media, electronic advertising via the internet, mobile phones or other electronic media have the benefit of being able to record data relating to the audiences' interaction with the advertisement. For example, internet based advertisements in the form of banner advertisements provide for the tracking of whether the advertisement was presented to the user as well as whether the user clicked on that advertisement to view more information about the subject of the advertisement. Advertising pricing models are often based on performance metrics relating to the number of views per advertisement."}
{"text": "Merchants find it useful to keep track of who their customers are, and how frequently their customers make purchases. Merchants regularly use this information to, for example, administer customer incentives programs, for marketing purposes, and to formulate advertising strategies.\nFor those customers that make purchases using a financial transaction card (e.g., a credit card, debit card, smart card, etc.), merchants potentially have available to them an easy way of obtaining both identifying and purchase frequency information. Every financial transaction card has stored on it data, such as the customer's personal account number (PAN), which is unique to a particular customer. Such financial transaction card data is extracted from the card and sent to the merchant during the course of a typical card-based transaction. Merchants could store that unique financial transaction card data, and use it to determine each time a purchase was made by a customer using a particular financial transaction card.\nHowever, merchants are currently unable to make use of financial transaction card data to keep track of customers and purchase frequency. For security reasons, current payment schemes for financial transaction cards have dictated rules that forbid merchants from storing financial transaction card data that contains information about a customer's account. There remains a need in the art for a simple method and system that would allow merchants to uniquely identify their customers using financial transaction card data, but that does not require that the merchants store that data."}
{"text": "1. Field of the Invention\nThe present invention relates in general to loadcarrying vehicles, and, more particularly, to a novel transformer caddie for transporting and locating transformers and other equipment.\n2. Description of the Prior Art\nModern underground residential distribution systems of electric power frequently include small pad-mounted distribution transformers located in unpaved areas behind homes, usually backyards, these areas being generally inaccessible to trucks and cranes typically employed in transporting and locating transformers and other relatively heavy electrical equipment. Periodic replacement of such transformers has, in the past, presented serious problems, particularly when replacement transformers must be moved through gates or other narrow passageways and into hilly terrain. Furthermore, the problems associated with transformer replacement are aggravated by mud and snow conditions.\nOne device used for transporting such distribution transformers is a small dolly or wheeled platform. However, since such transformers are frequently quite heavy, often weighing on the order of one ton or more, a rather large work crew has been found to be necessary to wheel and set a transformer into place on its pad using this device.\nAnother known device is a pole-carrying vehicle comprising an inverted U-shaped frame having a winch mounted on its top for suspending a load, the device having a pair of steerable front wheels similar to those found on a wagon. However, this device lacks maneuverability, particularly necessary for positioning a transformer onto its mounting pad. Furthermore, since this device is rather large, its lack of maneuverability prevents easy movement through narrow passageways."}
{"text": "(1) Technical Field\nThe invention relates to electronic circuits, and more particularly to bias circuits for cascode amplifiers, particularly radio frequency cascode amplifiers.\n(2) Background\nAmplifiers are a common component in radio frequency (RF) transmitters and receivers, and are frequently used for power amplification of transmitted RF signals and for low-noise amplification of received RF signals. For many RF systems, particularly those requiring portability (e.g., cellular telephones, WiFi-connected computers, cameras, and other devices, etc.), it has become common to use complementary metal-oxide semiconductor (CMOS) fabrication technology to create low cost, low power, complex integrated circuits (ICs). CMOS device technology improvements, such as silicon-on-sapphire (SOS) CMOS, silicon-on-insulator (SOI) CMOS, and ever-shrinking field-effect transistor (FET) device channel lengths, are putting the RF performance metrics of silicon-based CMOS transistors on par with rival gallium arsenide (GaAs) technologies.\nFIG. 1 is a simplified schematic diagram of a prior art amplifier circuit 100. The illustrated circuit 100 has a cascode common source architecture constructed from two series-connected FETs. An RF input signal, RFIN, is applied through a biasing and coupling circuit 102 to the gate of a FET M1, the drain of which is coupled to (and thus drives) the source of a second FET M2. The source of M1 is coupled thorough a degeneration inductor L to RF ground. The degeneration inductor L, which will also have an inherent resistance R (shown in dotted outline), performs several functions, including obtaining a good dynamic range (e.g., a good noise figure), achieving high sensitivity with low power consumption, helping make input impedance matching easier, and improving linearity.\nThe second FET M2 has its gate coupled to a voltage source Vgate (which may be VDD in many cases) and provides an amplified RF output signal, RFOUT, at its drain. In some applications, the RF input signal, RFIN, may be processed through an input impedance matching network 104 before being coupled to the gate of M1, and in some applications, the RF output signal, RFOUT, may be processed through an output impedance matching network 106. While the illustrated cascode circuit provides good isolation because there is no direct coupling from the output to the input, the circuit is not well suited to applications in which the DC supply voltage varies, such as applications where the DC supply voltage at the drain is actively modified to optimize operation at different power levels; examples are average power tracking (APT), envelope tracking (ET), and GSM power amplifier power control. In these examples, the supply voltage VDD may have a range that varies by a factor of 10 or more (e.g., from 4.5V to 0.4V, as one example).\nOne problem area for using silicon-based CMOS transistors in cascode amplifier circuits is creating bias circuits, particularly in cases in which the DC supply voltage varies and in the case of amplifiers with a transistor stack height of 3 or more. In general, a DC bias voltage applied to the gate of a FET sets the gate-source voltage VGS so as to provide a desired drain current. Silicon-based CMOS devices with short channel lengths provide fast RF response—a desirable characteristic in an amplifier—but suffer from poor output resistance characteristics due to the influence of the drain voltage on the gate of a transistor device. As such, they are not amenable to common open-loop bias techniques, such as current mirrors, without suffering from a large and undesirable mismatch between a reference device current and the output device current. In addition, the problem can be further exasperated in SOS and SOI CMOS technologies due to issues associated with the floating body or “kink” effect. In the floating body effect, the body of a transistor may form a capacitor with the insulated substrate; charge accumulation on this capacitor may cause the formation of parasitic transistors consuming unwanted parasitic currents and further degrading the output resistance of devices intended to be used as current sources.\nAnother challenge in designing bias circuits for silicon-based CMOS cascode amplifiers is to configure such circuits so that they are tolerant of supply and bias voltage variations, and bias current variations. Traditional cascode bias circuits include diode connected devices which do not accommodate a wide range of supply and bias voltages. Bias current variations may be driven by variations in process, temperature, and voltage (i.e., “PVT” effects), the end result of which is to cause fluctuations in bias current leading to fluctuations in RF performance. Another drawback to using silicon-based CMOS devices for amplifier applications, and particularly for power amplifier applications, is a relatively low breakdown voltage per device.\nYet another challenge in designing RF amplifiers is that, in many applications, the RF electrical environment of a transmitter, receiver, or transceiver is constantly changing. The changing characteristics of an RF signal path affects RF circuit performance metrics such as gain, linearity, noise figure, and power consumption. A system performance profile that ideally sets such metrics for one situation (such as high gain, moderate linearity, and low noise in the presence of a small received signal) can be completely inappropriate as the RF environment changes (such as with increasing received signal power). These situation transitions can occur many times during a single usage session of an RF circuit (for example, as a cellular phone moves relative to signal towers), and not adjusting to an appropriate performance profile as the RF environment changes can result in poor or non-functioning RF circuit performance.\nAccordingly, there is a need for bias circuits for silicon-based CMOS amplifier architectures that are tolerant of supply and bias voltage variations, bias current variations, and transistor stack height, and compensate for the poor output resistance characteristics of silicon-based CMOS devices with short channel lengths. There is also a need for amplifier architectures that can rapidly adapt to a changing RF electrical environment. Further, there is a need for improved silicon-based CMOS amplifier architectures having good isolation."}
{"text": "The orthogonal frequency-division multiplexing method (OFDM) has established itself in many areas for high-rate data transmission in radio systems with large bandwidth. It is used, for example, in digital audio broadcasting (DAB), digital video broadcasting-terrestrial (DVB-T), in wireless local area networks (WLAN) and 4G long term evolution LTE. The principle of OFDM consists of splitting the high-rate overall data stream into many low-rate data streams and transmitting them parallel via the corresponding number of orthogonal subcarriers. Compared to single-carrier methods, various advantages result for transmission via frequency selective multipath channels. In doing so, channel equalization can be implemented efficiently in the frequency range without using costly adaptive equalization filters. Furthermore, by introducing a guard interval (GI) between the OFDM symbols, inter-symbol interference (ISI) can be effectively prevented. Signal generation can be realized with the help of inverse discrete Fourier transform (IDFT) or its expense-favorable implementation, the inverse fast Fourier transform (IFFT).\nAn almost identical method is also used for line-tied transmission. It is known as “discrete multitone transmission” (DMT). It is used, for example, for broadband digital data transmission via subscriber connection lines (DSL). Because DMT can also be viewed as a form of OFDM, the term “OFDM” will be used for both module types, hereafter. Thus, the following remarks similarly refer to DMT systems.\nOne problem with the transmission via channels with distinct multi-path dispersion (e.g. by means of signals reflected at buildings during radio transmission or reflection onto lines during line-tied transmission) is that destructive interference can cause erasing of subcarriers (“spectral zeros”) in the channel). Because the symbols modulated onto the affected subcarriers cannot be detected correctly by the receiver, any longer, this causes a high symbol error rate (SER) or a high bit error rate (BER) in the data stream, which only slightly decreases with increasing signal to noise ratio (SNR). Thus, it cannot readily be compensated for by increasing the transmission capacity.\nAs a matter of principle, the problem can be avoided if the transmitter only modulates the data to the strong subcarriers and omits the weak or deleted subcarriers. This involves, however, that it knows the actual radio channel to the receiver (channel state information at transmitter's), which is not the case in many systems. The reasons for this are, for example, the higher complexity of the system connected therewith, and feedback information that may be used (transmission of channel information from receiver to transmitter), which cause a transmission overhead and have to be made available to transmitter with very little delay for a fast time-variable channel.\nGenerally, OFDM is coupled with forward error correction (FEC), which results in coded OFDM (COFDM). To do this, targeted redundancy is added to the sent data stream by means of suitable channel coding and used at receiver's to correct transmission errors. By means of FEC, BER can be significantly reduced, whereby the net data rate is decreased by the added redundancy—a defaulted net data rate involves correspondingly increased transmission resources. A linked FEC is oftentimes implemented in classic COFDM systems, where a convolution code is linked as inner code and a block code (e.g. Reed Solomon Code) is linked as outer code, for example. The inner code has the primary task to reduce the effects of the weak subcarriers or the spectral zeros. With the help of the outer code, the error rate is further reduced by several decimal powers.\nOne alternative to compensate for the influence of weak subcarriers is to spread data symbols prior to OFDM modulation in the frequency range. Every symbol is then no longer transmitted via a single but via all subcarriers, via which the spreading operation is being implemented. Any sub-carrier carries a linear combination of all transmission symbols within the block. Even in case of loss of several subcarriers, it is often times still possible to reconstruct the transmission symbols, and the BER strongly decreases around the operating point (as in the non-frequency selective case) with increasing SNR. This is referred to as diversity gain. The method is referred to as “code-spread OFDM” (CS-OFDM; see reference [1] —the stated references can be found in detail in Appendix 1 to this description), “spread OFDM” (SOFDM; see reference [2]) or “frequency domain spreading” (see reference [3]). The abbreviation CS-OFDM will be used, hereafter. Classically, M modulation symbols are spread onto/over N subcarriers, causing no loss in data rate. If M<N modulation symbols are spread onto N subcarriers, the system is referred to as “partially loaded” (PL) (PL CS-OFDM; see reference [1]). By means of partial load, further gains can be achieved at the expense of the data rate.\nCompared to COFDM systems, no redundancy has to be added to CS-OFDM systems to utilize frequency diversity. An additional FEC with the help of channel coding is, however, generally reasonable to further reduce the error rate. Because the channel code used in such case, does not have to compensate for the influence of the weak subcarriers any longer, the code rate can turn out larger for classic coding methods in connection with the spreading.\nCS-OFDM signals have unfavorable signal statistics, in common with OFDM signals. Interference of a plurality of subcarriers results in a poor ratio of instant performance to mean signal performance (signal statistics). This is often times indicated by the peak-to-average power ratio (PAPR), the ratio of maximum instant performance to mean signal performance or the crest factor (ratio of the maximum instant value to the root mean square value of the signal). The measurements take on high values for OFDM signals. Signal statistics deteriorates with increasing subcarrier number N.\nHigh requirements to the linearity of the components used in the system, and especially to the power amplifiers, result from poor signal statistics. Therefore, such components can often times only be operated in an inefficient operating point. One problem for the PL CS-OFDM is the fact that the PAPR—in comparison to the fully loaded system—can drastically continue to deteriorate.\nThe presentation of the problem in respect to a CS-OFDM system is the realization of an efficient, adaptive transmission by means of an implementation with low complexity:                high performance ability: good utilization of the frequency diversity of the frequency selective channel, low bit error rate with given SNR and given net data rate        good signal statistics (low PAPR): low requirements to the linearity of the analog hardware (above all at the power amplifiers) and performance specific operating point        high adaptivity: adaptivity possibility of the transmission (robustness, net data rate, redundancy) to the conditions of the transmission channel with low complexity and fine tuning        implementation with low complexity: low utilization of resources of digital signal processing and low requirements to the speed of digital signal processing, cost and performance efficient implementation        \nGenerally, the criteria cannot be optimized independently from another. Conventional technology knows several spreading methods with good properties, which will be described, hereafter.\nConventional technology describes Hadamard (also termed “Walsh-Hadamard” or “Hadamard-Walsh”; see references [1]-[4] and [5]), DFT (see references [3], [5]) as well as Vandermonde spreading (see reference [4]). These spreading methods are characterized by high performance ability. This applies at least to large spreading lengths, that is, if the spreading is being implemented by means of a high number of subcarriers. Modifications of the Hadamard transform and the discrete Fourier transform (DFT), the “rotated” Hadamard transform and the “rotated” DFT are suggested in reference [5] to improve performance ability for low spreading lengths.\nAn increase of CS-OFDM is introduced in references [1], [4]. It consists of feeding a low number of modulation symbols into the spreading operation, which are available as subcarriers, that is, in spreading M symbols onto N subcarriers, where M<N. The system is then referred to as “partially loaded” (PL CS-OFDM). The classic case, a fully loaded system, exists for M=N. Reference [1] suggests the use of PL CS-OFDM to realize further gains compared to OFDM and CS-OFDM, at the expense of the data rate.\nFor DFT spreading, recourse can be made to the fast Fourier transform (FFT), that is, an efficient algorithm for calculation of the discrete Fourier transform (DFT). The complexity amounts to O (N log2 (N)), with multiplications with the root of unity and additions that may be used. The Vandermonde spreading, a multiplication of a Vandermonde matrix of the dimension N×N with a symbol vector of the length N, can be realized with the complexity O (N log2 (N)) (multiplications and additions; see reference [17]). Hadamard spreading can be implemented with the help of the fast Walsh-Hadamard transform (FWHT). It involves N log (N) additions or subtractions and thus has the lowest computational complexity among the methods with full flexibility (fine-tuned partial load, selection of the carriers).\nAn implementation of the DFT spreading with particularly low complexity is described in reference [16]. In doing so, the multiple carrier system is reduced back to a single carrier system, so that the transmitter can be realized very easily, which, however, is not similarly applicable to the receiver in case a channel estimation and a channel equalization have to be implemented. Furthermore, the system is no longer fully flexible with respect to the partial load with full diversity and selection of the allocated subcarriers.\nFIG. 1 shows a block diagram by means of which signal processing is explained in a conventional CS-OFDM transmitter. By means of an introduction, it should be noted that the mathematical notation used in this description is explained in Appendix 2 to this description. The matrices and sequences, to which reference is being made, are explained, there, as well. For better understanding, cross references to the sections of Appendix 2 are indicated at several locations. The abbreviations used are listed in Appendix 3 to this description.\nFIG. 1 shows a processing chain 100 for generation of the spread OFDM signal. Mathematically, this can be illustrated with the help of vectors, matrices and corresponding operations. The stream of the (already coded, where applicable) data symbols ds is transformed with the help of a serial/parallel converter 102 in data vectors d of the length M. Spreading is being implemented at block 104. The type of spreading is entirely described by the spreading matrix D. Subsequently, the spread data vectors x are fed to IDFT 106. Blocks 104 and 106 can be combined in one block 108, which implements both the spreading and the IDFT. The combination of both operations is hereafter referred to as “overall transform”, and block 108, which implements such transform, as “overall transformer”. The transforming output vectors w can then be submitted to further operations in the digital baseband (e.g. inserting a Guard interval, fenestration). By means of a digital/analog converter, they are converted to analog signals and run through the typical processing chain of a transmitter for digital data processing, until emission via the antenna(s).\nMathematically, the spread data vector x results from a vector matrix multiplication of the spreading matrix D (dimension N×M) with the input vector d (length M, column vector):x=Dd. \nIDFT 106 can be illustrated as multiplication with the IDFT matrix F−1 (see section 2.10 of Appendix 2) of the dimension N×N. We have:w=F−1x and consequentlyw=F−1Dd=:Bd, with the multiplication of vector d with the matrix B=F−1D illustrating the overall transform 108.\nThe spreading methods known in the art (see, for example, references [1]-[5]) can be characterized by means of matrices, in particular Hadamard, Vandermonde, and DFT matrices, (see definitions in sections 2.8, 2.9 and 2.10 of Appendix 2). With a full load of the system (M=N), D is identical to the characterizing matrix. With a partial load (M<N), a partial block or sub-block of the characterizing N×N matrix is used for D, which block consists of the first M matrix columns (see references [1], [4]). Spreading methods, which are based on a spreading matrix, are referred to as matrix-based spreading methods, hereafter."}
{"text": "Often times a homeowner is in his/her vehicle driving and they cannot recall whether they shut the garage door. If the garage door is not shut while the owner is not home, the house may be easily broken into by a criminal.\nAdditionally, a homeowner may be in his/her home and ready to go to bed. However, the homeowner cannot remember if he/she shut the garage door. Again if the garage door is not shut, the homeowner's family may be subjected to intruders."}
{"text": "Earrings are not easily stored in any organized fashion because of their small size and odd shapes, and the fact that they are used in pairs. The usual storage is accomplished by simply keeping them loose in a box or drawer, with the result that it is difficult to find a matching pair.\nDisplay racks found in retails stores usually involve displaying small cards or boxes in which are fastened a single pair of earrings in a manner such that the earrings are not readily detachable. As such, these racks are totally unsuitable for use in the home where one wants to quickly find and remove from storage a selected pair of earrings. Insofar as is known, there is no such storage rack available at the present time, nor has one been disclosed in the prior art.\nIt is an object of this invention to provide an improved holder for earrings for use in the home. It is another object to provide a holder for earrings in which the earrings are easily pinned to a fabric. Still other objects will become apparent from the more detailed description which follows."}
{"text": "Interests in energy storage technologies have been increasingly higher recently. As applications are expanded to energy of mobile phones, camcorders and notebook PCs, and furthermore, to electric vehicles, efforts on the research and development of electrochemical devices have been more and more materialized.\nElectrochemical devices are fields receiving most attentions in such aspects and among these, development of secondary batteries capable of charge and discharge have been the focus of attention, and in developing such batteries, research and development on the design of new electrodes and batteries for enhancing capacity density and energy efficiency have been recently progressed.\nAmong currently used secondary batteries, lithium secondary batteries developed in early 1990s have received attentions with advantages of having high operating voltage and significantly higher energy density compared to conventional batteries such as Ni-MH, Ni—Cd and sulfuric acid-lead batteries using an aqueous liquid electrolyte.\nA lithium secondary battery has a structure of an electrode assembly including a positive electrode, a negative electrode and a separator interposed between the positive electrode and the negative electrode being laminated or wound, and is formed by embedding the electrode assembly in a battery case, and injecting a non-aqueous liquid electrolyte thereinto. The lithium secondary battery produces electric energy through oxidation and reduction reactions occurring when lithium ions are intercalated/deintercalated in the positive electrode and the negative electrode.\nIn a common lithium secondary battery, a negative electrode uses lithium metal, carbon and the like as an active material, and a positive electrode uses lithium oxides, transition metal oxides, metal chalcogen compounds, conductive polymers and the like as an active material.\nAmong these, a lithium secondary battery using lithium metal as a negative electrode mostly attaches lithium foil on a copper current collector, or uses a lithium metal sheet itself as an electrode. Lithium metal has low potential and high capacity, and has received much attention as a high capacity negative electrode material.\nWhen using lithium metal as a negative electrode, electron density non-uniformization may occur on the lithium metal surface when operating a battery due to various reasons. As a result, a branch-shaped lithium dendrite is produced on the electrode surface causing formation and growth of projections on the electrode surface, which makes the electrode surface very rough. Such lithium dendrite causes, together with battery performance decline, separator damages and battery short circuits in severe cases. As a result, a temperature in the battery increases causing a risk of battery explosion and fire.\nIn order to resolve such problems, researches such as introducing a polymer protective layer or an inorganic solid protective layer to a lithium metal layer, increasing a concentration of a salt of a liquid electrolyte, or using proper additives have been urgently required."}
{"text": "The present invention relates to an artificial cardiac valve.\nNumerous artificial cardiac valves employing flap valve devices of the ball, disc or flap type in the form of a hemi-disc, already exist.\nThis is the case, for example, of the valves described in Patent Applications EP 0 176 337 (CARBOMEDICS INC.), EP 0 023 797 (HEMEX INC.), EP 0 039 217 (MITRAL MEDICAL INTERNATIONAL, INC. and EP 0 050 971 (HEMEX INC.).\nFurthermore, GB 2 084 299 (HEMEX INC.) describes a cardiac valve comprising an annular element and a flap pivotally lodged inside said element as well as linking elements adapted to cooperate with guiding and restraining means provided on the internal lateral wall of the annular element.\nThese documents describe in particular valves of the type comprising a flap with single pivoting disc. The disc in question pivots by means of hinges in the form of protuberances disposed on the inner wall of the body of the valve and which, due to their configurations and arrangements, are capable of provoking phenomena of thrombosis.\nFurthermore, these valves are in a position of complete opening only when the flap has pivoted through a certain angle which is always less than 90.degree. about an axis of rotation which merges with its diameter. In fact, the limited angle of 90.degree. does not allow passive closure of the valve. Consequently, in the position of opening, half of the disc constituting the flap is engaged in the aorta, which may raise problems of flow and/or space requirements.\nMoreover, it is observed that these valves tend to oscillate due to the great distance between the axis of rotation and the most remote end of the flap."}
{"text": "1. Field of the Invention\nThe present invention relates to a polyvalent clamp chucks for machine-tools, i.e. a rotating part having retractable jaws able to rigidly hold a workpiece in a very precise position. Numerous types of chuck are already known, each being designed for a specific use.\n2. Discussion of Background Information\nA first type of chuck has a cylindrical body with two, three, four or six concentric gripping jaws guided radially in its front face. By their design, arrangement and large degree of movement, the jaws of these chucks can grip large workpieces. These jaw-type chucks are differentiated according to the type of drive means provided for radially actuating the jaws in response to axial displacement of a concentric internal control cylinder in the chuck body.\nA first kind of jaw-type chuck, the so-called wedge-type, as, for example, described in French patent publication FR-2 447 768, comprises a control cylinder terminating with a sleeve having, facing each jaw, an oblique outwardly-directed wedge in the form of a groove of inverted T cross-section, in which wedge is engaged the rear heel of a main jaw, also of T shape. As the main jaw is mobile only radially in the front face of the chuck body, it is easily understood that when the control cylinder is withdrawn, the lower part of the sleeve's wedge causes lowering of the jaw's correction heel. FR-2 431 338 describes a device for compensating for centrifugal force for this wedge-type of chuck, comprising a lever situated to the rear of and perpendicular to the jaw, with a short arm bearing in a recess in the rear face of the jaw and a longer arm that is pulled outwards by an inertia block under the action of centrifugal force.\nIn a second kind of jaw-type chuck, the so-called sliding block chuck, such as that described in European patent publication EP-0 215 350, the main jaw is engaged on one side (or on both sides) by an oblique tenon and mortice coupling with a part sliding parallel to the chuck axis, this block being actuated by a rear control cylinder also via a tenon and mortice coupling. To avoid jamming effects, the tenon and mortice coupling between the sliding block and the main jaw may have a trapezoidal cross-section. Alternatively, the oblique coupling of jaws/sliding block may also be provided by a series of parallel, oblique and interpenetrating toothings.\nAnother jaw-type chuck of the so-called lever type, such as that sold by Messrs GAMET under reference \"MX\", comprises a lever for each jaw, situated in the radial plane passing through the jaw, with a short arm engaged in a recess provided in the jaw's rear face, and a longer arm at right angles to the other, that engages in a recess provided in the external face of the control cylinder. Thus, by withdrawing the control cylinder, the short arms are caused to move down, thus bringing the jaws together. Advantageously, one or all of the levers further comprises, in extension of their short arm and beyond the pivot, an inertia block that enhances gripping of the jaws under the action of centrifugal force.\nA second general type of chuck of the so-called pincer-clamp type, such as that described in European patent publication EP-0 258 771, comprises several pincers in the form of segments of a hollow truncated cone, whose external surfaces come to bear in an also-conical fixed part of the chuck body. The ends of these pincers are provided with hooks by means of which they can be pulled by the control cylinder towards the inside of the fixed conical part of the chuck. This type of chuck is more particularly designed for machining workpieces from bars or parts of medium dimensions requiring uniform gripping over a large surface.\nAnother general type of chuck, the so-called \"front-driving\" type, comprises a coaxial body whose end is fitted centrally with a centering point and, adjacent to this point, a series of radial driving knives. The centering point may be fixed with the mobile driving knives driven by a rearwardly-located jack, or vice versa. Cooperating with a similar facing body, these parts are designed to grip by its lateral faces a workpiece to be machined over its entire external surface. Frequently, these chucks are mounted on a jaw-type chuck, either by being directly fixed against the flat front face by bolts, or by being gripped in the jaws themselves. However, the centering of this fitted jaw may create some difficulties, which may slow down the installation operation.\nAnother general type of chuck, the so called expansible type, comprises a body fitted against the chuck's front face. An external part of this body can expand, or an internal part can contract, under the action of a displaceable control body or under the action of hydraulic pressure, so as to seize a workpiece to be machined on its external surface.\nA first type of expansible mechanical chuck comprises a simple external cone forming an integral part of the basic body and over which a slit part is pushed or pulled, according to the inclination of the cone. The converse arrangement can also be envisaged, namely a fixed pincer with an inner cone that is pushed or pulled.\nA second type of expansible mechanical chuck of the so-called \"double cone\" type comprises a pincer gripped between a conical part of a main body and a second mobile conical part situated at the end of a pull-piece. By pulling the pull-piece, thus bringing the two conical parts towards one another, it is possible to radially displace the pincers so they can grip from the inside a hollow part to be machined on its outer surface.\nA third type of expansible mechanical chuck comprises, also in a fitted body, an elastic membrane deformed by hydraulic pressure produced by displacement of a piston in a cylinder integrated in the fitted body, which piston is manually or automatically actuated by a jack.\nFrom the above, it can be understood that a mechanic's workshop having to machine numerous small series of different parts, as, for instance, is usual for subcontractor's workshops, must imperatively be equipped with several types of chuck, which requires a large investment.\nMoreover, changing from one type of chuck to another is often long and intricate, and increases the dead time between production runs. Furthermore, it would be possible to allow entirely automatic machines to run continuously, even when the operator is absent during breaks, the workpieces being loaded and unloaded by the machine as long as no manual adjustment is required when passing from machining with orthogonal jaws to machining with oblique jaws, or vice versa.\nLastly, upon each change of the clamping chuck, the position of the gripping jaws on the main jaws must accordingly be set to a procedure that may vary from one chuck manufacturer to another."}
{"text": "1. Field of the Invention\nThe present invention relates to methods and systems for organizing and displaying financial information.\n2. Description of the Related Art\nFidelity Investments' Web site provides a visual representation of stocks in the stock market, which is branded as “Map of the Markets”. Each of the stocks is represented by a rectangular shape. An aggregate of the rectangular shapes is also rectangular. Though all of the stocks are represented in a compact form, it is unclear how a stock represented by one rectangle specifically relates to stocks represented by neighboring rectangles.\nPCQuote.com introduced a graphical representation of a single equity, branded as “Sniper”. The representation has a bulls eye design comprised of 13 sectors, each sector representing an area of reconnaissance. The areas of reconnaissance include: stock flow which indicates a relationship between a number of shares that traded on the bid side versus a number of shares traded on the offer side, quadrant distribution which indicates a total volume of shares traded distributed across four quadrants to show what prices are attracting the most activity, price rotation which numerically indicates when price rotation activity favors buyer or sellers, relative price rotation which visually indicates when price rotation activity favors buyers or sellers, time rotation which measures the acceptance of higher or lower prices, relative time rotation which shows a current time period bias, “to target” which indicates when a stock price is within ½ point of a forecast high or low price, “to exit” which indicates a most extreme stop-out price, buy/sell imbalance which measures where a majority of trades are occurring, look-out signal which indicates a day's direction for an equity, gatekeeper which provides a logical stop-out price, stock and market direction which points out a path of least resistance, and “center of the scope” which indicates whether buyers or sellers are dominant across the board."}
{"text": "1. Field of the Invention\nThis invention relates to a prevention of buildup condensate and mold on the inner surface of the beehives caused by temperature and humidity conditions. In particularly, the invention relates to a beehive climate and moisture control system which prevents freezing air from entering the beehive and balances the temperature and humidity surrounding the cluster.\nMore particularly, this invention relates to a creation of balancing relationship between the sources of the heat, specially constructed winter ventilation inner cover, ventilation, ability to use the subject of invention at remote apiaries without permanent sources of electrical power.\n2. Background Art\nAs the temperature inside the bee hive reach between 10 and 12° C. bees start to establish cluster. Because there is no air circulation inside the hive, temperatures at different levels of the hive would have tendency to stay different. Cluster can be built, based on bees aspiration toward warmer areas, at higher levels (elevations). As soon as cluster is set up, it is instinctively moving upward. The frames with food, located under the cluster will never be reached by bees, which would resulted by the end or even middle of cold season to lack of food and colony starve to death for that reason.\nExcessive moisture within a beehive creates an unhealthy environment for bees during winter time. High humidity promotes decay of the wooden structures, mold growth and may contribute to dysentery and other health problems. During cold periods the bees form a cluster within and around the combs of honey, which supply food when foraging is no longer possible. When bees consume honey from comb frames, they heat up bee cluster only, but not surrounding area. Although the bees produce heat, the temperature inside beehive is essentially at the same temperature as the outside. Bees from the core feed on stored honey and gradually work to the outer layer of the cluster to change places with cold bees. While the bees can survive dry cold alone, a cold wet bee cannot survive.\nTo protect, bees from excess humidity condensation derived from accumulated moisture drops needs to be prevented in a way that they can drop onto semi-donnant bees. The moisture is accumulated from the heat of metabolism of the living cluster as they feed on stored honey. The small amount of heat generated nevertheless causes vapor to rise and condense upon the first cold surface it contacts.\nWhen air with high water vapor content and relatively high Per Dew point met the surface of the ceiling, walls as well as frames surfaces with lower Per Due Dew point, it is start to condensate on their colder, practically equal to outside temperature surfaces. Large amount of honey, pollen which is organic substances and perfect food source in the bee hive create ideal environment for mold grow. There is a need to keep constantly removing any excessive moister out of the bee hive, but at the same time as much as possible avoiding to cool down the interior of the bee hive, especially cluster surrounding area. That why the main sensing device is not temperature sensor but humidity sensor, which is actually through the level of RH in used air watching two main parameters of humidity and temperature inside the beehive.\nA need exists in the current art for an improved beehive climate and moisture control system that may be used during winter time. The present invention solves the problem by a beehive system according to claims 1 to 18."}
{"text": "The present invention relates to a portable terminal, an information processing apparatus, a content display system, and a content display method.\nCellular phones capable of browsing internet sites and television (TV) programs have come into wide use today. On the other hand, there have been increasingly employed television sets capable of connecting to networks. In this environment, a technique in which a content received by a cellular phone is viewed by a television set with higher usability and a technique in which a cellular phone is employed as a remote control device to operate a television set are under discussion.\nFor example, JP-A-2006-286855 describes a cellular phone including a broadcast receiving function. The cellular phone transmits a history of viewed programs and playback indication signals thereof to a video playback device. According to the cellular phone, a program which the user viewed outside his or her house can be again viewed on a large screen after returning home. It is also possible that when the user viewed part a program outside home during a free period of time, the user can view its subsequent part, parts before and after the viewed part, or the entire program from the start point to the end point hereof after returning home by using a recording and reproducing device in the house.\nJP-A-2005-135346 describes a technique for use in a system including a cellular phone, an information processing apparatus, and a display. When mail data including Uniform Resource Locator (URL) information is received, the cellular phone transfers the mail data via radio communication to the information processing apparatus. On receiving the mail data, the information processing apparatus displays the mail data on the display. If an operation is conducted for the URL information in the mail data, an operation is conducted to access a radio communication network other than a network used for the radio communication on the basis of the URL information, and resultantly obtained information is displayed on the display.\nJP-A-2006-338406 describes a technique in which a communication terminal device including another function is ordinary employed as a remote control device to remotely control other communication terminal devices.\nAdditionally, JP-A-2007-74265 describes a videoconference system using a cellular phone connected via the internet to a videoconference server to communicate information with the server and a display device capable of, when connected via a communication function to the cellular phone, displaying the information. The videoconference system controls devices connected thereto via the videoconference server and a network. The videoconference system includes an information display module to display on a first display of the cellular phone a first videophone screen on which the information is displayed using letters. The information display module also displays on a second display of the display device a second videophone screen on which image information included in the information communicated as above is displayed. JP-A-2007-74265 describes a system in which a monitor camera is remotely controlled by a cellular phone, and a video image produced from the camera is displayed on a display screen of a Personal Computer (PC)"}
{"text": "1. Field of the Invention\nThe present invention relates to an image signal processing apparatus which is used for an electronic still camera or the like, and in which the transfer control of an image signal sent from an image pickup element is improved in connection with the encoding of the image signal.\n2. Description of the Prior Art\nAs an apparatus for converting an optical image into an image signal through an image pickup element and for storing the image signal as digital data into a storage device, there is, for example, a digital-type electronic still camera. In the electronic still camera, image data compression is essential in order to store a bulky amount of digital image data in a storage device of a given capacity. As the method of image data compression, conventionally, there have been proposed various methods. For example, known methods are a DPCM system using a correlation between picture elements adjacent to each other, a DCT system in which a picture area is divided into n x n picture element areas (blocks) composed of n columns of picture elements and rows of picture elements and a correlation in which every block is used.\nFIG. 9 is a block diagram showing an example of the signal processing in an electronic still camera employing the DPCM system. This example will be described hereunder.\nAn image signal produced from an image pickup element 11 such as a CCD or the like is amplified by an amplifier 12, and the amplified signal is separated into three color signals, for example, red, green, and blue (R, G, B) color signals, by a color separation circuit 13. Then, the color signals R, G, and B are processed by gamma circuits 14 and white balance circuits 15 respectively, converted into digital signals by A/D conversion circuits 16 respectively, and stored in a field memory 17. The field memory 17 is used as a buffer memory for a DPCM encoding circuit 18.\nThe DPCM encoding circuit 18 serves to encode the image data which is converted into the digital signal and read from the field memory 17 in the order of color frames or in the order of color lines. FIG. 10 shows the configuration of the DPCM encoding circuit 18. That is, in the DPCM encoding circuit 18, a predicted value X.sub.0 (obtained by predicting a value at the present point of time on the basis of a sampled value/values at a sample point/points before the present point of time) is subtracted from a sampled value x.sub.0 at the present point of time by means of a subtracter 81 to thereby obtain a prediction error signal .epsilon..sub.0. The prediction error signal .epsilon..sub.0 obtained as a result of the subtraction is outputted after being non-linearly converted by a non-linear converter 82. The predicted value may be obtained, for example, by a method using a sampled value at a sampling point one picture element before the present point of time (hereinafter referred to as a previous-sample-prediction method).\nReference numeral 83 designates a local decoder portion which serves to prevent accumulation of errors generated in the decoding operation. In the local decoder portion 83, a non-linear reverse converter 84 serves to reversely convert an output of the non-linear converter 82, and an adder 85 serves to add a prediction error signal which is an output of the non-linear reverse converter 84 to the above-mentioned predicted value, the result of addition being supplied to a predictor 86.\nThe image data encoded in the DPCM encoding circuit 18 is written, through a buffer memory 19, into an image data storage medium such as an IC memory card or the like in a storage device 20. The buffer memory 19 serves to finally write the image data into a storage medium 20a while controlling the rate of data transfer between the DPCM encoding circuit 18 and the image data storage medium 20.\nIn FIG. 9, a CCD driving circuit 21 serves to operate a CCD used as the image pickup element 11. Further, a timing generation circuit 22 is controlled by a system control circuit 23 so as to supply timing signals at predetermined points of time to various circuits 13 through 21 to thereby operate those circuits 13 through 21.\nIn such a configuration, the gamma circuit 14, the white balance circuit 15, and the A/D conversion circuit 16 are required for each of the color signals of R, G, and B. Further, all the color signal R, G, and B separated by the color separation circuit 13 are simultaneously outputted from the color separation circuit 13. Accordingly, if the field memory 17 were not provided, it would be necessary to provide the DPCM encoding circuit 18 for each of the color signals of R, G, and B. It is however difficult to actually realize such a configuration, and, therefore, it is necessarily required to provide, in a camera, a buffer memory, such as a field memory 17 as mentioned above or a frame memory, of a large capacity.\nThus, in the conventional case, it has been necessary to provide three sets of the gamma circuits 14, the white balance circuits 15, and the A/D conversion circuits 16 for the respective R, G, and B color signals and the field memory 17 or the like which is a large-capacity buffer memory. As a result, the circuit scale and cost is prevented from being reduced so that the apparatus cannot be made inexpensive."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of leakage coaxial cables which are used for radio communication in closed in areas such as the interior of buildings, tunnels or underground markets, where no ordinary radio waves can be received.\n2. Description of the Prior Art\nA leakage coaxial cable typically has slots formed in the exterior thereof which are spaced at predetermined intervals along the conductor, so that electromagnetic waves propagating inside the coaxial cable are partially radiated from the conductor into an external space through the slots.\nWhen a fixed signal source is connected to the leakage coaxial cable, the signal is radiated into the external space and may be received by a mobile station running near the leakage coaxial cable. In addition, a signal transmitted by the mobile station may be received by the fixed station through the leakage coaxial cable.\nA common application for leakage coaxial cables has been in disaster warning and prevention system for use in building, tunnels and underground markets. It is important in such applications to make certain that the cables are fireproof. The ability of prior art leakage coaxial cables to withstand high heat has been limited. Consequently, there has existed a long and unfilled need in the prior art for a method of making leakage coaxial cables that have transmission characteristics which will not degrade during a fire related emergency.\nA conventional leakage coaxial cable has an external conductor with slots formed therein to radiate an electromagnetic wave which is propagating inside the cable outwardly into an external space. The external conductor is disposed coaxially around an internal conductor with an insulating member therebetween. The external conductor is covered by a protective sheath. In order to minimize the transmission loss of the electromagnetic wave, the insulating member is preferably made of a low-loss plastic material, such as polyethylene or polystyrene. The external conductor is preferably made of high-conductivity material such as alminum or copper. A polyester film is laminated on the external conductor with adhesive to compensate for the decrease of the mechanical strength of the conductor which is caused by the formation of the slots. The protective sheath is preferably made of polyethylene or polyvinyl chloride.\nIf a leakage coaxial cable thus constructed encounters a fire, the protective sheath will burn away and the external conductor will be directly exposed to the flames. The polyester will burn, and the plastic insulating member will melt. The molten plastic will flow through the slots formed in the external conductor, ignite and drop from the cable while burning. The burning molten plastic may actually contribute to the spread of the fire and may burn the skins or clothes of persons fighting the fire or running away from the fire.\nOne example of a prior art method of making a leakage coaxial cable refractory has been disclosed in Japanese Utility Model Application Publication No. 16682/1977. In that method, a heat-resistant tape which is made of an inorganic material such as asbestos is spirally wound between a polyethylene insulating member and an external conductor. The heat-resistant tape prevents the melting of the polyethylene insulating material for a long time. In addition, the heat-resistant tape maintains the insulation between the internal and external conductors even after the polyethylene insulating member has been molten. Therefore, the radio communication properties of the cable can be maintained unchanged for some time after the occurrence of the fire.\nHowever, the leakage coaxial cable disclosed in Japanese Utility Model Application Publication No. 16682/1977 is still not an ideal solution to the problem discussed above, since the polyethylene insulating material will eventually ooze out through the stitch lines of the non-organic tape or the seams of the spirally wound tape and flow out through the slots. In addition, since the tape is relatively thick, on the order of 0.25 to 0.5 mm, the dielectric value between the internal and external conductors is large, causing large transmission losses.\nAnother example of a prior art conventional leakage coaxial cable is disclosed in Japanese Utility Model Application (OPI) No. 3537/1980. In that cable, a heat-resisting tape of polyimide resin is spirally wound between an external conductor and a polyethylene insulating member. Because of the presence of the heat-resisting tape, the internal and external conductors are not short-circuited even if the polyethylene insulating member is molten.\nHowever, that leakage coaxial cable also has the problem of the polyethylene insulating member oozing out through the seams of the spirally wound tape, flowing out through the slots and dropping from the cable while burning."}
{"text": "1. Field of the Invention\nThe present invention is broadly concerned with improved breakaway sign posts which can be used to support various types of signs in roadway environments. More particularly, the invention is concerned with such sign posts, as well as the coupling assemblies used to interconnect upper and lower sign post sections, wherein a magnetic coupling arrangement is employed in lieu of complex and costly mechanical couplers.\n2. Description of the Prior Art\nPosts of many different varieties are used in association with roadways. Such posts may be employed to support signs, providing information to vehicle operators, and barriers, directing vehicles away from hazardous areas. Most posts along a roadway are fixed to the ground by digging a post hole, inserting the base of the post into the hole and then filling the hole with concrete or other material to support the post.\nSuch roadway posts are often inadvertently hit by vehicles traversing the roadways, which usually result in breaking of the posts. When a post breaks, the top portion thereof can rotate toward the vehicle and possibly enter the vehicle's passenger compartment. This can result in serious or even fatal injuries to the occupants.\nA further problem with roadway posts currently in use is the difficulty in replacing the post once it has been broken. To reinstall a damaged post, the original concrete must be dug up and removed. Then, a new post is inserted into the hole and fresh concrete is poured to fix it in place. This is labor intensive and requires several man-hours of labor for each broken post.\nBreakaway posts have been devised in the past for use in supporting signs, while providing a breakaway function during a sudden vehicle impact. For example, U.S. Pat. No. 6,264,162 describes a successful design of this character. This breakaway post assembly includes a specialized breakaway collar having strategically located lines of weakness which yields during a vehicle impact. Additional breakaway post designs are described in the following U.S. Pat. Nos. 5,480,121, 3,349,531, 3,820,960, 3,846,030, 4,490,062, 4,926,592, 4,928,446, 5,088,683, 5,160,111, 5,214,886, 5,484,217, 5,535,555, 5,782,040, and 5,988,598.\nThese prior art breakaway posts are all characterized by mechanical connection structure between upper and lower separate post sections. While providing a useful breakaway function, once impacted by a vehicle, it is necessary to reinstall the collar or related structure to put the post back into an operative condition. Moreover, in certain of the prior designs, it is impossible to repair the breakaway post, and a new installation is necessary. In other instances, the prior breakaway posts do not break away in any predetermined fashion or direction, i.e., it is possible that the broken away portion may fall toward the vehicle rather than away from it.\nThere is accordingly a real and unsatisfied need in the art for an improved breakaway post assembly which avoids complex mechanical collar arrangements and permits easy repair of a post after it has broken away, and which is designed to cause the broken away post section to fall in a predetermined direction away from the impacting vehicle."}
{"text": "Modem semiconductor integrated circuits usually involve multiple layers separated by dielectric (insulating) layers, such as of silicon dioxide or silica, often referred to simply as an oxide layer, although other materials are being considered for the dielectric. The layers are electrically interconnected by holes penetrating the intervening oxide layer which contact some underlying conductive feature. After the holes are etched, they are filled with a metal, such as aluminum, to electrically connect the bottom layer with the top layer. The generic structure is referred to as a plug. If the underlying layer is silicon or polysilicon, the plug is a contact. If the underlying layer is a metal, the plug is a via.\nPlugs have presented an increasingly difficult problem as integrated circuits are formed with an increasing density of circuit elements because the feature sizes have continued to shrink. The thickness of the oxide layer seems to be constrained to the neighborhood of 1 .mu.m, while the diameter of the plug is being reduced from the neighborhood of 0.25 .mu.m or 0.35 .mu.m to 0.18 .mu.m and below. As a result, the aspect ratios of the plugs, that is, the ratio of their depth to their minimum lateral dimension, is being pushed to 5:1 and above.\nFilling such a hole with a metal presents two major difficulties.\nThe first difficulty is filling such a hole without forming an included void, at least with a filling process that is fast enough to be economical and at a low enough temperature that doesn't degrade previously formed layers. Any included void decreases the conductivity through the plug and introduces a substantial reliability problem. Chemical vapor deposition (CVD) is well known to be capable of filling such narrow holes with a metal, but CVD is considered to be too slow for a complete filling process. Physical vapor deposition (PVD), alternatively called sputtering, is the preferred filling process because of its fast deposition rates. However, PVD does not inherently conformally coat a deep and narrow hole. A fundamental approach for applying PVD to deep holes is to sputter the metal on a hot substrate so that the deposited material naturally flows into the narrow and deep feature. This process is typically referred to as reflow. However, high-temperature reflow results in a high thermal budget, and in general the thermal budget needs to be minimized for complex integrated circuits. Further, even at high temperatures, the metal does not always easily flow into a very narrow aperture.\nThe second difficulty is that an aluminum contact is not really compatible with the underlying semiconductive silicon. At moderately high temperatures, such as those required for the reflow of aluminum into the narrow hole, aluminum tends to diffuse into the silicon and to severely degrade its semiconductive characteristics. Accordingly, a diffusion barrier needs to be placed between the aluminum and the underlying silicon.\nThese problems are well known and have been addressed by Xu et al. in U.S. patent application, Ser. No. 08/628,835, filed Apr. 5, 1996, incorporated herein by reference in its entirety, which is a continuation in part of U.S. patent application, Ser. No. 08/511,825, filed Aug. 7, 1995 now U.S. Pat. No. 5,962,923.\nAs shown in the cross-sectional view of FIG. 1, a contact hole 10 having an aspect ratio defined by its depth 12 and its width 14 is etched through a dielectric layer 16 to an underlying substrate 18, which in the more difficult situation includes a surface layer of silicon. In the hole filling process, the contact hole 10 is conformally coated with a titanium (Ti) layer 20, a titanium nitride (TiN) layer 22, and a graded (TIN.sub.x) layer 24, that is, the graded layer 24 begins at its bottom as TiN but its top portion is nearly pure Ti. These three layers form a tri-layer barrier 26, which provides both the conformality and the adhesion to the underlying layers, as well as sufficient wetting for the after deposited aluminum. A Ti layer 20, after siliciding at the sufficiently high annealing temperature, forms a good ohmic contact with the underlying silicon substrate 18. Thereafter, a metal layer 28 is sputter deposited into the hole 10 so as to fill it without voids. That is, the tri-layer barrier 26 sufficiently wets to the after filled aluminum that it readily flows into the hole 10 at a moderate temperature while the tri-layer barrier 26 nonetheless provide a sufficient diffusion barrier between the aluminum 28 and the underlying silicon 18.\nAccording to Xu et al., the wetting quality of the three layers 20, 22, 24 is enhanced by depositing them in a high-density PVD reactor. On the other hand, they recommend that the aluminum layer 28 be sputter deposited in a conventional PVD chamber with a low plasma density. In particular, they recommend that the aluminum layer 28 be deposited as two layers in an improved two-step cold/warm version of a conventional sputtering process. In the first cold step, a seed layer 30 of aluminum is sputter deposited at a substrate temperature below 200.degree. C. so as to conformally coat the underlying barrier tri-layer 26 with a fairly uniform aluminum layer. In the second warm step, a filling layer 32 of aluminum is sputter deposited at higher temperatures so as to reflow and fill the contact hole 10. An advantage of the tri-layer barrier 26 grown by ionized metal plating (IMP) is that the warm Al filling layer 32 can be filled at temperatures below 400.degree. C., even as low as 350.degree. C. according to the reported data. The warm layer 32 can be deposited at a fairly high rate so as to improve the system throughput. Because the two aluminum layers 30, 32 differ primarily in their different deposition temperatures, they are likely deposited within a single conventional PVD chamber capable only of developing a low-density plasma. Also, the two deposition can be performed continuously, with the temperature being ramped up during the deposition. As a result, the two Al layers 30, 32 have no clear boundary between them.\nIn the context of contact hole filling, a high-density plasma is defined in one sense as one substantially filling the entire volume it is in and having an average ion density of greater than 10.sup.11 cm.sup.-3 in the principal part of the plasma. The conventional plasma-enhanced PVD reactor produces a plasma of significantly lower ion density. Although high-density plasmas are available in a number of different types of reactors, they are preferably obtained in an inductively coupled plasma reactor, such as the type shown in schematical cross section in FIG. 2. For reasons to be described shortly, this is referred to an ionized metal plasma or plating (IMP) reactor.\nAs shown in this figure, which is meant only to be schematical, a vacuum chamber 40 is defined principally by a chamber wall 42 and a target backing plate 44. A PVD target 46 is attached to the target backing plate 44 and has a composition comprising at least part of the material being sputter deposited. For the deposition of both titanium (Ti) and titanium nitride (TiN), the target 46 is made of titanium. A substrate 48 being sputter deposited with a layer of a PVD film is supported on a pedestal electrode 50 in opposition to the target 46. Processing gas is supplied to the chamber 40 from gas sources 52, 54 as metered by respective mass flow controllers 56, 58, and a vacuum pump system 60 maintains the chamber 40 at the desired low pressure.\nAn inductive coil 62 is wrapped around the space between the target 46 and the pedestal 50. Three independent power supplies are used in this type of inductively coupled sputtering chamber. A DC power supply 64 negatively biases the target 46 with respect to the pedestal 50. An RF power source 66 supplies electrical power in the megahertz range to the inductive coil 62. The DC voltage applied between the target 46 and the substrate 48 causes the processing gas supplied to the chamber to discharge and form a plasma. The RF coil power inductively coupled into the chamber 40 by the coil 62 increases the density of the plasma, that is, increases the density of ionized particles. Magnets 68 disposed behind the target 46 significantly increase the density of the plasma adjacent to the target 46 in order to increase the sputtering efficiency. Another RF power source 70 applies electrical power in the frequency range of 100 kHz to a few megahertz to the pedestal 50 in order to bias it with respect to the plasma.\nArgon from the gas source 54 is the principal sputtering gas. It ionizes in the plasma, and its positively charged ions are attracted to the negatively biased target 46 with enough energy that the ions sputter particles from the target 46, that is, target atoms or multi-atom particles are dislodged from the target. The sputtered particles travel primarily along ballistic paths, and some of them strike the substrate 48 to deposit upon the substrate as a film of the target material. If the target 46 is titanium or a titanium alloy and assuming no further reactions, a titanium film is thus sputter deposited, or in the case of an aluminum target, an aluminum film is formed.\nFor the sputter deposition of TiN in a process called reactive sputtering, gaseous nitrogen is also supplied into the chamber 40 from the gas source 52 along with the argon. The nitrogen chemically reacts with the surface layer of titanium being deposited on the substrate to form titanium nitride.\nAs Xu et al. describe in the cited patent application, a high-density plasma, primarily caused by the high amount of coil power applied to the chamber 40, increases the fraction of the sputter species that become ionized as they traverse the plasma, hence the term ionized metal plating. The wafer bias power applied to the pedestal 50 causes the pedestal 50 to become DC biased with respect to the plasma, the voltage drop occurring in the plasma sheath adjacent to the substrate 48. Thus, the bias power provides a tool to control the energy and directionality of the sputter species striking the substrate 48.\nXu et al. disclose that the Ti/TiN/TiN.sub.x barrier tri-layer 26 should be deposited in an ionized metal process plating (IMP) process in which the various power levels are set to produce a high-density plasma. They observe that an IMP barrier tri-layer 26 as shown in FIG. 1, when deposited in the contact hole 10, promotes the reflow of aluminum into the contact hole 10 when the aluminum is subsequently deposited in a conventional PVD reactor, that is, one not using inductively coupled RF power and not producing a high-density plasma. This superior reflow is believed to require two characteristics in a narrow aperture. The barrier layer needs to adhere well to the underlying SiO.sub.2 or Si so as to form a continuous, very thin film. The aluminum needs to wet well to the barrier layer so that it flows over the barrier layer at relatively low temperatures.\nAlthough the TiN IMP barrier tri-layer offers significant advantages in promoting reflow of subsequently deposited conventional PVD aluminum, as processing requirements become even more demanding, further improvement of reflow into narrow apertures is desired."}
{"text": "The present invention relates to a lamp tube support structure and, more particularly, to a lamp tube support piece of low cost, simple structure, and easy assembly.\nAs shown in FIGS. 1 and 2, a prior art office business machine like a scanner la has a lamp seat 11a disposed therein. A lamp tube 12a used as a light source for scanning documents is disposed on the lamp seat 11a. The lamp tube 12a can be divided into single-tube type and double-tube type. A rack 13a for supporting the lamp tube 12a is disposed below the lamp tube 12a. In order to meet the compactness requirement of present products, the lamp tube 12a must be designed to be thinner. Therefore, the lamp tube 12a easily breaks because of vibration of the scanner 1a or careless fall during the transportation process.\nAccordingly, the above lamp tube fixing structure has inconvenience and drawbacks in practical use. The present invention aims to resolve the problems in the prior art.\nThe primary object of the present invention is to provide a lamp tube support structure, which is slipped onto a lamp tube to add support positions at two ends of the lamp tube for effectively preventing breakage of the lamp tube. Moreover, clamping effect can be achieved to protect the lamp tube from impact from all directions.\nTo achieve the above object, the present invention provides a lamp tube support structure, which comprises mainly a support piece applying to office business machines for supporting and fixing a lamp tube thereof. The lamp tube support structure comprises a first half and a second half. The first half has a gap for placing a lamp tube. A clamping portion is formed at two extension ends of the gap. Grooves are disposed at left and right ends of the second half. A wing portion is formed below the grooves. After the wing portion is bent, it can be jammed into a hole of a rack for supporting the lamp tube so as to be erectly fixed on the rack."}
{"text": "This application is based upon and claims the benefit of priority from the prior Japanese Patent Applications No. 2000-182564, filed Jun. 19, 2000; No. 2000-182566, filed Jun. 19, 2000; and No. 2001-056268, filed Mar. 1, 2001, the entire contents of all of which are incorporated herein by reference.\nThe present invention relates to an onboard X-ray CT apparatus which can be mounted in a motor vehicle, trailer, or container, a container for mounting the X-ray CT apparatus, and a motor vehicle for mounting the X-ray CT apparatus.\nX-ray CT apparatuses for obtaining the tomographic image of an object to be examined by exposing the object from his/her surroundings to X-rays, measuring the amount of X-rays having passed through the object, and performing image reconstruction processing for the obtained projection data by using a computer are effective for detection of a disease in its early stage and determination of an appropriate therapeutic plan. Further, a so-called X-ray CT apparatus onboard vehicle, e.g., a motor vehicle such as a bus which contains an X-ray CT apparatus is moved to a disaster-stricken area where image diagnosis for the presence/absence of a fracture or an injury to a viscus is done. Such a motor vehicle is also exploited for itinerant diagnosis or group examination in a depopulated area such as a mountainous area.\nRecently, the development of helical scanning X-ray CT apparatuses enables photographing in a pectoral region within one breath. Along with this, strong demands have arisen for implementation of an X-ray CT apparatus onboard vehicle which contains a helical scanning X-ray CT apparatus for group examination in order to detect a pectoral disease in its early stage.\nA known example of a conventional X-ray CT apparatus onboard vehicle is disclosed in Utility Model Registration No. 3,027,769, and schematically shown in FIG. 1. A gantry 2a and bed 2b of an X-ray CT apparatus are mounted at the center of a large-bus type vehicle 1 in the longitudinal direction. A power supply device 3 and high-voltage generator 4 are installed in the space between a driver\"\"s seat 18 at the front portion of the vehicle and the gantry 2a and bed 2b. An X-ray monitor 5 and console 6 are installed at the back portion of the vehicle 1. A doorway 12 and lift device 15 for an object who uses a stretcher or wheelchair are formed in the back wall surface of the vehicle 1. A doorway 13 for an object who can walk by himself/herself is also formed in the backward side wall surface of the vehicle 1.\nOther X-ray CT apparatus onboard vehicles are disclosed in Jpn. Pat. Appln. KOKAI Publication Nos. 6-154202 and 8-127282.\nIn the X-ray CT apparatus onboard vehicle disclosed in Utility Model Registration No. 3,027,769, the power supply device 3 is installed in the same room as the gantry 2a and bed 2b of the X-ray CT apparatus, which impairs safety and gives annoying noise to an object during an examination.\nThe X-ray CT apparatus disclosed in Jpn. Pat. Appln. KOKAI Publication No. 6-154202 performs X-ray CT photographing while an object sits on a chair, and is not suitable for X-ray CT photographing of a severely wounded object. The X-ray CT apparatus disclosed in Jpn. Pat. Appln. KOKAI Publication No. 8-127282 is an industrial X-ray CT apparatus mounted in a large trailer, and is not suited for itinerate diagnosis or group examination in a depopulated area such as a mountainous area owing to the road situation.\nThese conventional X-ray CT apparatus onboard vehicles do not consider any protection means for the movable portions of the mounted X-ray CT apparatuses, and the movable portions of the X-ray CT apparatuses may be damaged by vibrations or shocks during traveling of the X-ray CT apparatus onboard vehicle.\nIn the conventional X-ray CT apparatus onboard vehicles, a CT operation room and the like are designed in accordance with a mounted X-ray CT apparatus, which makes it difficult to exchange the X-ray CT apparatus with another type apparatus upon upgrading the apparatus.\nIt is an object of the present invention to provide an X-ray CT apparatus suitably mounted in a motor vehicle, trailer, or container.\nAccording to the present invention, there is provided an X-ray CT apparatus which can be mounted in a motor vehicle, trailer, or container, comprising locking mechanism unit for locking at least one of movable portions of the X-ray CT apparatus, detectors for detecting locking states of the locking mechanism unit, and a notifying unit for representing the states of the locking mechanism unit on the basis of outputs from the detectors.\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out hereinafter."}
{"text": "Bibliographic details of the publications referred to by the author in this specification are collected alphabetically at the end of the description.\nThe reference in this specification to any prior publication (or information derived from it), or to any matter which is known, is not, and should not be taken as an acknowledgment or admission or any form of suggestion that that prior publication (or information derived from it) or known matter forms part of the common general knowledge in the field of endeavour to which this specification relates.\nMalignant tumours, or cancers, grow in an uncontrolled manner, invade normal tissues, and often metastasize and grow at sites distant from the tissue of origin. In general, cancers are derived from one or only a few normal cells that have undergone a poorly understood process called malignant transformation. Cancers can arise from almost any tissue in the body. Those derived from epithelial cells, called carcinomas, are the most common kinds of cancers. Sarcomas are malignant tumours of mesenchymal tissues, arising from cells such as fibroblasts, muscle cells, and fat cells. Solid malignant tumours of lymphoid tissues are called lymphomas, and marrow and blood-borne malignant tumours of lymphocytes and other hematopoietic cells are called leukaemias.\nCancer is one of the three leading causes of death in industrialised nations. As treatments for infectious diseases and the prevention of cardiovascular disease continues to improve, and the average life expectancy increases, cancer is likely to become the most common fatal disease in these countries. Therefore, successfully treating cancer requires that all the malignant cells be removed or destroyed without killing the patient. An ideal way to achieve this would be to induce an immune response against the tumour that would discriminate between the cells of the tumour and their normal cellular counterparts. However, immunological approaches to the treatment of cancer have been attempted for over a century with unsustainable results.\nAccordingly, current methods of treating cancer continue to follow the long used protocol of surgical excision (if possible) followed by radiotherapy and/or chemotherapy, if necessary. The success rate of this rather crude form of treatment is extremely variable but generally decreases significantly as the tumour becomes more advanced and metastasises. Further, these treatments are associated with severe side effects including disfigurement and scarring from surgery (e.g. mastectomy or limb amputation), severe nausea and vomiting from chemotherapy, and most significantly, the damage to normal tissues such as the hair follicles, gut and bone marrow which is induced as a result of the relatively non-specific targeting mechanism of the toxic drugs which form part of most cancer treatments and is a major limiting factor for dosage\nStill further, common chemotherapy drugs do not significantly penetrate into tissue further than about 70 microns from the blood supply (Primeau et al. Clin. Canc. Res. 2005, 11:8782-8788; Minchinton et al. Nat. Rev. Cancer 2006, 6:583-592). The rapid growth and poor vascular development of most solid tumours puts many tumour cells well beyond the capacity of the drugs to penetrate the tissue. Critically, many cells experience sub-lethal doses, allowing them to survive and to develop drug resistance.\nSolid tumours cause the greatest number of deaths from cancer and mainly comprise tumours of the linings of the bronchial tree and the alimentary tract that are known as carcinomas. In the year 2000 in Australia, cancer accounted for 30% of male deaths and 25% of female deaths (Cancer in Australia 2000, 2003) and it accounted for 24% of male and 22% of female deaths in the US in year 2001 (Arias et al. 2003, National Vital Statistics Reports 52:111-115). Solid tumours are not usually curable once they have spread or ‘metastasised’ throughout the body. The prognosis of metastatic solid tumours has improved only marginally in the last 50 years. The best chance for the cure of a solid tumour remains in the use of local treatments such as surgery and/or radiotherapy when the solid tumour is localised to its originating lining and has not spread either to the lymph nodes that drain the tumour or elsewhere. Nonetheless, even at this early stage, and particularly if the tumour has spread to the draining lymph nodes, microscopic deposits of cancer known as micrometastases may have already spread throughout the body and will subsequently lead to the death of the patient. In this sense, cancer is a systemic disease that requires systemically administered treatments. Of the patients who receive surgery and/or radiotherapy as definitive local treatment for their primary tumour and who have micrometastases, a minor proportion may be cured or at least achieve a durable remission from cancer by the addition of adjuvant systemic treatments such as cytotoxic chemotherapy or hormones.\nConventionally, solid cancer has been treated locally with surgery and/or radiotherapy, and during its metastatic stage with systemically administered cytotoxic drugs, which often interfere with the cell cycle of both normal and malignant cells. The relative selectivity of this approach for the treatment of malignant tissues is based to some extent on the more rapid recovery of normal tissues from cytotoxic drug damage. More recently, the targeted therapy of cancer has aimed to improve the therapeutic ratio of cancer treatment by enhancing its specificity and/or precision of delivery to malignant tissues while minimising adverse consequences to normal non-malignant tissues. Two of the major classes of targeted therapy are (i) the small molecule inhibitors such as the tyrosine kinase inhibitors imatinib mesylate (Glivec®), gefitinib (Iressa®) and erlotinib (Tarceva®), and (ii) the monoclonal antibodies (mAb) such as rituximab (Mabthera®) and trastuzumab (Herceptin®).\nIn parallel to the development of targeted therapies, combining at least two conventional anti-cancer treatments such as chemotherapy and radiotherapy in novel ways has been another approach to the development of cancer therapeutics. By exploiting synergistic interactions between the different modalities of treatment, combined modality treatment seeks to improve treatment efficacy so that the therapeutic ratio for the combined treatment is superior to that for each of the individual treatments.\nCombined modality treatment using external beam radiation and radiosensitising chemotherapeutic drugs such as 5-fluorouracil and cisplatin (chemoradiotherapy) has improved survival in a number of solid tumours such as those of head and neck, lung, oesophagus, stomach, pancreas and rectum because of both improved local tumour control and reduced rates of distant failure (TS Lawrence. Oncology (Huntington) 17:23-28, 2003). Although radiosensitising drugs increase tumour response, they also increase toxicity to adjacent normal tissues, which is especially true of the potent new generation radiosensitisers, gemcitabine and docetaxel. However, decreasing the radiation volume allows cytotoxic doses of gemcitabine to be better tolerated clinically (Lawrence 2003, supra). Chemoradiotherapy may overcome mutually reinforcing resistance mechanisms, which may only manifest in vivo.\nRadioimmunotherapy (RIT) is a systemic treatment that takes advantage of the specificity and avidity of the antigen-antibody interaction to deliver lethal doses of radiation to cells that bear the target antigen. Radio-isotopes that emit β-particles (e.g. 131Iodine, 90Yttrium, 188Rhenium, and 67Copper) are usually used to label monoclonal antibodies (mAb) for therapeutic applications. The energy from □-radiation is released at relatively low intensity over distances measured in millimeters (Waldmann, Science 252:1657-1662, 1991; Bender et al., Cancer Research 52:121-126, 1992; O'Donoghue et al. Journal of Nuclear Medicine 36:1902-1909, 1995; Griffiths et al. International Journal of Cancer 81:985-992, 1999). Thus, high-energy □-emitters such as 90Yttrium are useful for the treatment of larger and heterogeneous solid tumours (Liu et al. Bioconjugate Chemistry 12:7-34, 2001). Research interest in radioimmunotherapy has been reawakened because in spite of the low radiation doses delivered, significant and unexpected biological effects of RIT upon surrounding host cells have been observed (Xue et al. Proceedings of the National Academy of Sciences of the United States of America 99:13765-13770, 2002). Furthermore, the lower but biologically effective dose of radiation delivered by RIT had greater cytocidal effects than a larger dose of radiation conveyed as external beam radiotherapy (Dadachova et al., PNAS 101:14865-14870, 2004). Nonetheless, the efficiency of RIT as a treatment for solid tumours may be hampered by the low penetration of antibody through the tissue barriers that surround the target antigen in the tumour, which will consequently extend circulatory half life of the antibody (Britz-Cunningham et al. Journal of Nuclear Medicine 44:1945-1961, 2003). Furthermore, RIT is often impeded by the heterogeneity of the target antigen's expression within the tumour. Thus, although RIT affords molecular targeting of tumour cells, the major limitation of RIT remains the toxicity that may result from large doses of radiation that are delivered systemically in order to achieve sufficient targeting (Britz-Cunningham et al. 2003, supra; Christiansen et al. Molecular Cancer Therapy 3:1493-1501, 2004). Altogether, a useful therapeutic index using RIT has proven difficult to achieve clinically (Sellers et al. Journal of Clinical Investigation 104:1655-1661, 1999).\nTumour associated antigens, which would allow differential targeting of tumours, while sparing normal cells, have also been the focus of cancer research. Although abundant ubiquitous antigens may provide a more concentrated and accessible target for RIT, studies adopting this have been extremely limited.\nThe development of nanoparticle technology was also hailed as an exciting new frontier in terms of the development of new and effective cancer treatments. However, although previous attempts at using particulate material, such as nanoparticles, to target tumours for either diagnostic or therapeutic purposes have been extensive, in the context of therapeutics there has, disappointingly, been minimal success. With diagnostics, relatively shallow penetration of the particles into the tumour has been sufficient to achieve the objective of visualising the tumour. However, in terms of the delivery of a therapeutic agent, such shallow penetration has not been sufficient to effectively deliver the agent throughout the tumour, in particular to the interior of the tumour, as is required if total tumour destruction is to be achieved. In relation to therapeutics, specifically, conjugation of particles to a wide variety of different materials has so far failed to live up to the promise of achieving effective tumour penetration, this being an essential prerequisite for a therapeutic to have any chance of effectiveness.\nSignificant effort has also been made to take advantage of the enhanced permeability and retention (EPR) effect of tumours as a means to develop an effective therapeutic. Without limiting the present invention to any one theory or mode of action, this is a well described phenomenon based on the notion that certain sizes of molecules, typically liposomes or macromolecular drugs, tend to preferentially accumulate in tumour tissue. The general explanation for this phenomenon is that, in order for tumour cells to grow quickly, they must stimulate the production of blood vessels. VEGF and other growth factors are involved in cancer angiogenesis. Tumour cell aggregates of sizes as small as 150-200 μm become dependent on blood supply carried by neovasculature for their nutritional and oxygen supply. These newly formed tumour vessels are usually abnormal in form and architecture. They comprise poorly-aligned defective endothelial cells with wide fenestrations, lacking a smooth muscle layer, or innervation with a wider lumen, and impaired functional receptors for angiotensin II. Furthermore, tumour tissues usually lack effective lymphatic drainage. All these factors will lead to abnormal molecular and fluid transport dynamics, especially for macromolecular drugs. Accordingly, it has been thought that one way to achieve selective drug targeting to solid tumours is to exploit these abnormalities of tumour vasculature in terms of active and selective delivery of anticancer drugs to tumour tissues, notably defining the EPR effect of macromolecular drugs in solid tumours. Due to their large molecular size, nanosized macromolecular anticancer drugs administered intravenously escape renal clearance. Often they cannot penetrate the tight endothelial junctions of normal blood vessels, but they can extravasate in tumour vasculature and become trapped in the tumour vicinity. Nevertheless, the EPR effect has not been efficiently or successfully harnessed.\nVarious nanoparticles have been designed which are directed to achieving efficient cellular endocytosis. However, even if this is achievable, the issue of tissue penetration is still a separate one which, to date, has not been successfully overcome. The general notion of the use of a nanoparticle as a vector for delivery of a drug is widely discussed in the literature but, in the absence of achieving deep tumour penetration, is of limited value.\nEven where effective tumour distribution of a drug is achieved (by whatever means) a further problem has been the fact that neoplastic cells within solid tumours can exhibit a slowed metabolism. This means that even if a cytotoxic drug penetrates to these cells, if it is not effectively metabolised it will have a limited impact on the viability of the tumour.\nAccordingly, there is an urgent and ongoing need to develop improved systemic therapies for solid cancers, in particular metastatic cancers.\nIn work leading up to the present invention it has been determined that particulate material which is maintained in a dispersed state by a stabiliser is able to achieve deeper penetration into solid tumour models than has previously been achievable using nanoparticle technology. This has enabled the development of an effective means for treating solid tumours, both primary and metastatic, based on the co-administration of a cellular toxin with the particulate material. By either sequentially or simultaneously delivering this toxin, deeper penetration and therefore more extensive cellular exposure to the toxin is achieved. By virtue of the less effective reticuloendothelial clearance which is associated with tumours, a form of targeted treatment is effectively achieved. Still further, it has been observed that the toxin uptake by tumours penetrated by the particles of the present invention is effective, suggesting upregulation of tumour cell metabolism. Accordingly, the method of the present invention provides a means for achieving a more effective localised delivery and uptake of a cellular toxin to a tumour and its metastases in a manner which is characterised by significantly improved outcomes and/or reduced side effects relative to those which would normally be expected in the context of conventional treatment of an equivalent type of tumour. This is an extremely significant development since current protocols directed to treating metastatic disease are based on the non-targeted systemic delivery of chemotherapeutic agents."}
{"text": "1. Field of the Invention\nThis invention relates to a clutch cover assembly for an automotive vehicle or the like, and more particularly to a clutch cover assembly which includes a diaphragm spring.\n2. Description of the Related Art\nA conventional clutch cover assembly is disclosed in FIG. 8. This conventional art shown in Japanese Utility Model Laid Open Print No. 6256832, includes a cover 101 connected to an input delve member 102, a pressure plate 103 arranged in and connected to the cover 101, a clutch disk 104 disposed between the pressure plate 103 and the input drive member 102 and connected to an output driven member (not shown), and a diaphragm spring 105 disposed between the cover 101 and the pressure plate 103 wherein the pressure plate 103 is biased into the direction of the clutch disk 104.\nHowever, the diaphragm spring 105 is radially positioned by contacting with an inner wall surface 101a of the cover 101. Therefore, an outer peripheral surface 105a of the diaphragm spring 105 must be cut so as to position a center thereof accurately. Thus, the positioning of the center of the diaphragm spring 105 is very difficult.\nAlso, when a facing 107 of the clutch disk 104 is worn or damaged, the pressure plate 103 is moved in the direction of the clutch disk 104 by the spring force of the diaphragm spring 105. The diaphragm spring 105 is also moved in the direction of the clutch disk 104 by the spring force thereof. Thereby, an inner peripheral portion 105b of the diaphragm spring 105 contacts the clutch disk 104 and presses it against the input drive member 102 even when the facing 107 is worn away. Thus, the clutch disk 104 is damaged. In order to avoid this problem, a stopper 108 which is provided on an outer peripheral portion 101b of the cover 101 prevents the movement of the pressure plate 103 in the direction of the clutch disk 104 by more than a certain amount. However, the stopper 108 requires a large radial space."}
{"text": "Enterprise interworking typically involves the establishment of end-to-end connectivity among sites of an enterprise (or site) over a “wired” infrastructure. This type of connectivity has been traditionally enabled through various routing protocols, such as open shortest path first (OSPF), routing information protocol (RIP), border gateway protocol (BGP), and the like. Depending on the types of underlying networks, the operation of these routing protocols may be different but, in all cases, each protocol seeks to achieve the propagation of reachability information among interconnected sites so as to establish direct end-to-end connectivity between the sites. For security purposes, enterprises also seek to establish this direct end-to-end connectivity through private networking services. Typically, this capability is offered through addressing and routing techniques, such as multiprotocol label switching, virtual private networking (MPLS/VPN) services, which essentially enable service providers to logically partition their networks—using, for example, virtual routing and forwarding (VRF)—to permit sites with overlapping internet protocol (IP) addressing to share the same physical networking infrastructures, but maintain privacy through logical networking virtualization techniques.\nAs wireless networking capabilities continue evolve and offer increased reliability and throughput, enterprises are increasingly considering wireless connectivity as a potential source of both primary and secondary access to their enterprise (or private) networks. As such, suitable wireless interconnecting networks must be able to support the aforementioned direct end-to-end connectivity, as well as enable private networking services. Unfortunately, network virtualization and the ability to support multiple enterprises having overlapping or identical IP addressing spaces have been limited to the “wired” domains.\nTherefore, there is a need for an approach that extends private enterprise networking to wireless interconnecting domains."}
{"text": "1. Field of the Invention\nThis invention relates to a camera with a flashlight emitting device.\n2. Description of the Prior Art\nIn recent years, cameras of the lens shutter type having a built-in flash device are on the market which include a zoom lens having a focal length which can be varied continuously. In a camera of the type mentioned, a flash coverage or irradiation angle of flashlight from the flash device is normally varied in response to a zooming operation of the zoom lens.\nHowever, a camera wherein a flash coverage is varied in response to a zooming operation while a mode in which flashlight is ot emitted is selected by a user not only consumes power wastefully but also may make the user feel disagreeable. Besides, such camera is inconvenient to a user in that, upon use thereof, it can photograph an object only within a zooming range of the zoom lens.\nFurther, in a camera of the lens shutter type having a built-in flash device and including a zoom lens, after a flash coverage is varied in response to a zooming operation, a focusing operation of the lens is performed in response to completion of such zooming operation.\nHowever, if a focusing operation is held inhibited while a flash coverage is being varied in response to a zooming operation, then so much time is required until initiation of photographing that a shutter chance may possibly be lost.\nMeanwhile, in a camera of another type which has a built-in flash device and is capable of mounting an interchangeable lens thereon, a flash coverage of the flash device is fixed to a predetermined angle. Since such flash coverage is fixed in the camera of the type just mentioned, a certain focal length of the lens may result in unnatural exposure.\nOn the other hand, various cameras are also on the market which include an interchangeable camera and which control, when a flash device is mounted, the flash coverage of the flash device to a flash coverage corresponding to a focal length of the interchangeable lens in response to operation of a release button.\nIn a camera of the type mentioned just above, a change of the focal length of the interchangeable lens will not cause a variation of the flash coverage, and the flash coverage is not varied until the release button is operated. Accordingly, when it is intended to operate the release button suddenly to take a photograph, such photographing may be effected while the flash coverage remains in such a condition wherein it does not correspond to a current focal length of the lens. Such photographing may result in irregular flash coverage.\nAlso, cameras are on the market which have a built-in lens having a variable focal length and a built-in flash device having a variable flashlight irradiating range and wherein the focal length of the lens and the flash coverage of the flash device are varied in response to operation of a release button or of a focal length varying operating member such as a power zooming button.\nIn a camera of the type just mentioned, the lens and the flash device are connected mechanically to each other such that variation of the focal length and variation of the flashlight irradiation range may be performed simultaneously by a single electromagnetic driving means. Such mechanism can be realized with a camera which has a built-in lens and a built-in flash device. In the case of a camera which includes an interchangeable lens, however, it is very difficult to mechanically connect a lens and a flash device to each other particularly because different interchangeable lenses are interchangeably mounted on the camera."}
{"text": "Conventionally, dressings for the treatment of wounds or pressure sores are essentially flat dressings which are sufficiently conformable to be applied to flat or curved areas of the body and possess a sufficient absorbent capacity to absorb wound exudate. Such flat dressings are not very suitable for applying on protruding parts of the body such as elbows, heels or the tips of fingers or toes.\nPublished EP application No. 0 573 708 discloses a process for the forming of means comprising hydrocolloid adhesives for wound treatment. This reference discloses production of means or dressings for treatment of wounds which means or dressings can be prepared in a continuous manner.\nPublished EP application No. 0 092 999 discloses an occlusive wound dressing which may be used for treatment of wounds on the skin and a particulate or granular material being able to act together with the exudate of the wound. Wounds secreting vast amounts of liquid may be treated by firstly applying the granulate material and then to cover the wound with the material for treating wounds. This reference is silent with respect to three-dimensional dressings and concentrate on wounds or damages from which a considerable amount of liquid is secreted.\nU.S. Pat. No. 4,876,748 discloses a dressing containing a water-soluble or water-swellable hydrocolloid, a water-insoluble, viscous elastomeric binder and optionally a tackifier resin which is bevelled along all outer edges and optionally also, if annular, along the inner edge.\nU.S. Pat. No. 4,367,732 discloses a skin barrier consisting of an elastic film, secured to the film a layer of at least weakly elastic adhesive material, these two components together having low resistance to quick deformation and rapid recovery to substantially the original shape after deformation, the plastic properties of the adhesive thereby being compensated by the elasticity of the film. Optionally there may be a protective cover on the other side of the adhesive layer. The adhesive material consists of a hydrocolloid dispersed in a continuous phase consisting of a mixture of a physically cross-linked elastomer, a hydrocarbon resin tackifier, a plasticizer for the elastomer, an antioxidant and optionally an oily extender. This reference does not disclose three-dimensional dressings.\nDK patent No. 147,226 discloses a dressing having an adhesive layer and a water-proof polymer film. The dressing comprises a layer of pressure-sensitive adhesive and a water-impermeable flexible polymer film in which the adhesive comprises a rubbery elastomer in which is distributed a water-soluble or water-swellable hydrocolloid or a mixture of hydrocolloids, a tackifying resin and a plasticizer. Between the adhesive layer and the flexible polymer film is a layer of flexible foam having semi-open cells. The reference does not disclose three-dimensional dressings.\nPublished EP application No. 0 055 023 discloses an antiseptic adhesive comprising a natural or synthetic viscous rubbery mixture having one or more water-swellable hydrocolloids and antiseptic agents. It has been shown that having one or more antiseptic agents incorporated in the adhesive, the probability of bacterial growth will be reduced considerably and in many cases be totally avoided. The healing effect does not seem to be adversely afflicted. This reference does not mention three-dimensional dressings.\nU.S. Pat. No. 3,348,541 discloses a prepared bandage complementary to the end of a finger including at least the first joint thereof, and having an adhesive strip for holding the bandage on the finger.\nU.S. Pat. No. 2,847,005 discloses a surgical dressing for forming a finger cot. The dressing comprises an elongated flexible member permanently formed as an arcuate trough for disposition longitudinally along the outer surface of a finger, said trough including an open rear end and a closed forward end portion, and comprising a part to be folded over the finger once placed in the trough. The part to be folded over the finger comprises a strip of adhesive along its edge adhering to the finger.\nU.S. Pat. No. 2,875,785 discloses a finger tip bandage having a central cornpress or protecting pad with adjoining adhesive areas in the form of a flat sheet having two pairs of notches to allow bending the bandage to fit the finger.\nU.S. Pat. No. 2,243,422 discloses a finger bandaging unit comprising a gauze pad having substantially rectangular body and a tab portion extending from one edge thereof and adhesive-surfaced securing means carried by said pad.\nPublished International Patent Application No. WO 92/05756 discloses a conformable wound dressing of concavo-convex or cup-like shape comprising a body facing layer of an apertured elastomeric material having a concave nonadherent wound contacting surface, an outer layer of bacteria impermeable moisture vapor transmitting elastomeric film having an outwardly facing convex surface and an extensible intermediate layer of absorbent material which may be a sheet of polymer foam. The dressing is made from a flat blank.\nNone of the above references disclose nor indicate the existence of three-dimensional dressings for covering a protruding part of the body, said dressing having a three-dimensional part having an adhesive surface for contacting the wound.\nOne object of the invention is to provide a dressing, for example a finger tip or toe tip dressing allowing the choice of materials and form of the dressing ensuring an improved healing and a dressing being simple to produce and to apply to the wound to be covered. Another object of the invention is to provide a dressing which may prevent e.g. wearing or abrasion damages, e.g. on heels or elbows, and are provided with a surface which may be adapted to the environment in which the dressing is to be used giving a longer effective time of use for the dressing between the change the dressing. A further object of the invention is to provide a dressing which comprises emollients or e.g. retinoids for treating or preventing formation of psoriasis, eczema, callous skin, corns or blisters. A still further object of the invention is to provide processes for the preparation of such dressings."}
{"text": "1. Field of the Invention\nThis invention relates to the purification of brominated aromatic compounds such as decabromodiphenyl ether and pentabromophenol and more particularly to a process for purifying aromatic compounds containing residual free bromine and by-product hydrogen bromide.\n2. Description of the Prior Art\nHigh levels of purity are required for brominated aromatic compounds which have found utility as flame retardant agents in polymer compositions. In particular, it is important that such brominated products have extremely low levels of residual impurities such as free bromine, hydrogen bromide, retained catalysts, by-product bromine containing derivatives and the like since the presence of such impurities can have undesirable effects on the polymer compositions in which such agents are used. Purity is particularly important from the standpoint of color, and thermal stability under the processing conditions in commercial molding operations is likewise an essential parameter.\nHigh purity is an especially important consideration in the case of decabromodiphenyl ether, a flame retardant agent that has found wide application as an additive for high impact polystyrene used in television cabinets and other consumer appliances.\nThe preparation of decabromodiphenyl ether by the direct bromination of diphenyl oxide using excess bromine itself as a reaction solvent, as disclosed in British patent specification No. 1,411,524, published Oct. 29, 1975, and in Stepniczka, U.S. Pat. No. 3,965,197, issued June 22, 1976, produces highly brominated derivatives containing undesirably large amounts of occluded free bromine and by-product hydrogen bromide.\nThe cited British specification describes the use of various chemical purification treatments (e.g., injection of sulfur dioxide to convert occluded bromine into hydrobromic acid and injection of ethylene to convert residual bromine into dibromoethane). The cited U.S. patent merely describes the recovery of crude brominated product and gross separation of bromine therefrom by injecting superheated steam and by washing with hot dilute hydrochloric acid and hot water. The patent does not describe purification of the crude reaction project.\nTraditional purification methods such as recrystallization techniques are usable only with difficulty with materials such as decabromodiphenyl ether because its limited solubility in available solvents makes recrystallization both cumbersome and uneconomical.\nBurk, U.S. Pat. No. 3,733,366, granted May 15, 1973, describes a procedure for decolorizing brominated biphenyl obtained by treating biphenyl with bromine or bromine chloride in methylene chloride solvent in the presence of aluminum chloride catalyst. The patentee describes heating the product to a temperature between about 100.degree. C. and about 160.degree. C. with a preferred temperature range being about 110.degree. C. to about 150.degree. C. (Col. 2 lines 7-13). The disclosed heating step may take place optionally in the presence of solvents such as ethylene dibromide, toluene and xylenes. The patent is concerned only with decolorizing brominated biphenyl and contains no disclosure that would aid one in purifying materials such as decabromodiphenyl ether using a heat treating step.\nBritton, U.S. Pat. No. 2,022,634, issued Nov. 26, 1935 relates to the halogenation of diphenyl oxide and describes (page 2, col. 1, lines 3-7) heating the reaction mixture to above 75.degree. C., preferably between 250.degree. and 350.degree. C., to destroy bromine addition products. Britton's products are not completely brominated, nor does the patentee recognize that free bromine, rather than bromine addition products, may be removed by such a heating step.\nNagy et al, U.S. Pat. No. 3,752,856, issued Aug. 14, 1973, is directed to a process for producing brominated compounds in which an intimate mixture between bromine and the organic starting material is obtained by physical mixing with a sigma type blender. After bromination, the apparatus is flushed with dry air, preferably under lowered pressure, to remove the residual hydrobromic acid and, if applicable, excess bromine, and the whole mass is allowed to cool while continuing the grinding operation. The brominated aromatic product is said to be removed in the form of a powder as a rule. (Column 3, lines 5-11).\nThe patentee suggests that gaseous ammonia be passed into the apparatus to neutralize \"the hydrobromic acid and possibly the bromine retained by the product\" and goes on to describe purification of \"[t]he crude product\" by \"washing with acidified water or better still, by wet grinding in the presence of a dilute inorganic acid, followed by washing with water and by drying.\" The reference also suggests that \"[r]ecrystallization in an appropriate solvent offers another possibility of purification.\"\nNagy et al's Example III is a preparation of decabromodiphenyl in which the crude product was \"heated at the temperature of 200.degree. C. under normal pressure and under a nitrogen flow . . . ,\" the resulting product containing over 2 weight percent ammonium and aluminum bromide impurities. In Example IX of the patent, crude decabromodiphenyl was first subjected to dry air at 70.degree. C. under vacuum (20 mm Hg) followed by heating at 150.degree.-200.degree. C. \"under a flow of air.\" Recovery was completed by washing the \"crude product\" with aqueous NaOH. In Example IV, which is directed to decabromodiphenyl ether, purification is carried out by recrystallization in chlorobenezene.\nBrackenridge, U.S. Pat. No. 3,833,674, issued Sept. 3, 1974, describes a process for polybromination of aromatic compounds including diphenyl ether by reacting an aromatic compound with bromine in the presence of methylene bromide solvent. The patent states (column 2, lines 11-15) that \"a main improvement\" of the invention is the isolation of the product by precipitation from the reaction mixture by adding methanol or the like as a precipitant.\nBrackenridge describes four approaches to enhancement of product purity. (Column 4, line 46, et seq.) More particularly, the patent suggests the exclusion of light from the methylene bromide and methanol and from the reaction and recovery steps, minimization of temperature, and use of distilled solvent and precipitating agent. (Id., lines 60-68).\nAs a result, the prior art has failed totally to appreciate the applicability of the techniques of the present invention to the purification of brominated aromatic compounds such as decabromodiphenyl ether. In particular, the prior art has failed to recognize the importance of the crude decabromodiphenyl ether being thermally stable, much less that thermal stability is directly controlled by the substantial exclusion of aliphatic and alicyclic hydrocarbon impurities from the raw materials employed. The art further fails to disclose the importance of precise particle size control of the brominated aromatic product to be subjected to the heating step in achieving desired elimination of bromine, bromide, and other impurities.\nAccordingly, it is a primary object of this invention to obtain a process for producing purified brominated aromatic compounds that is superior to the techniques that heretofore have been employed.\nAnother object is to provide a method of the character described that may be economically employed in purifying decabromodiphenyl ether.\nA still further object is to provide a method of obtaining thermally stable decabromodiphenyl ether that may be purified in accordance with the present invention."}
{"text": "This invention relates to an electronic taximeter, and more particularly to an electronic taximeter which provides simple changes in tariff display conditions in various operating modes (as determined by an actuated tariff key and a prior actuated tariff key).\nTariff display conditions in an electronic taximeter are generally fixed with regard to various operating modes of the taximeter. For instance, \"PAID\" is displayed upon actuation of a paid key, \"20% UP\" upon actuation of a 20% up key and \"OCCUPIED\" upon actuation of an occupied key. There may be a need to change these tariff display conditions in the respective operating modes of the taximeter due to tariff system reforms. In this case, a problem persists with the prior art taximeter in that the circuit arrangement thereof must itself be reconstructed whenever such changes are necessary."}
{"text": "In modern motor vehicles, systems for assisting a parking process, in brief parking assistance systems, are often used. Conventionally, in this case, surroundings of a vehicle are checked for the presence of a suitable parking space by sensor-based measurement while passing a parking space, and the vehicle is parked in a selected parking space using the parking assistance system according to a decision by the driver. There are parking assistance systems that can carry out the entire parking process automatically. With autonomously driven vehicles, automatic parking assistance systems are also occurring more frequently.\nTypically, in an observation phase, a parking assistance system of a vehicle determines a suitable parking space in which the vehicle fits with the dimensions thereof, wherein adjacent regions are included within a determination that are sufficient for boarding and exiting and loading, and that guarantee a safety distance from other vehicles. If a suitable parking space is found, in a calculation phase, a trajectory is calculated in which a steering angle and vehicle speeds within defined sub segments of the path to be traversed are incorporated. Depending on the prevailing dimensions of the vehicle and the parking space, the trajectory can be divided into sub segments in which the vehicle may move back and forth multiple times. Furthermore, correction movements can be incorporated into the trajectory. In a parking phase, the parking assistance system starts a movement of the vehicle that is controlled according to the calculated trajectory into the parking space. Conventional parking assistance systems operate during this phase according to the open-loop principle (open control loop), with which, once calculated, the trajectory is not changed when movement of the vehicle has been started. If a vehicle bounding the parking space moves during the parking process, the parking process cannot be dynamically adjusted to the changed parking space. The object of the present disclosure is to improve conventional methods based on the use of parking assistance systems."}
{"text": "p-Diethylbenzene (p-DEB) is used as the desorbent in the Parex process (Derosset, et al., Sep. Sci. Technol. 15, 637 (1980)) for the separation of p-Xylene for xylene isomers and as the cross-linking agent in the synthesis of polymeric resins. One of the ways to produce p-DEB is to separate mixed DEB supplied from ethylbenzene alkylation process (Wang, et al., U.S. Pat. No. 4,950,835). Because of the very close boiling points possessed by DEB isomers: for example, ortho-diethylbenzene (o-DEB) is 183.degree. C.; meta-diethylbenzene (m-DEB) is 180.degree.-181.degree. C.; and para-diethylbenzene (p-DEB) is 184.degree. C., it is difficult to have a complete separation by distillation. The cryogenic separation is an alternative, however, it is in general costly.\nAt present, while several adsorption processes have been developed for separation of diethylbenzene isomers, the selective adsorption process using zeolite as a adsorbent in liquid phase is most commonly used and it is recognized as a most economic one. For example, disclosed in the U.S. Pat. No. 4,051,192 (1977), Neuzil, et al. used about 70 cm.sup.3 of the zeolite X ion-exchanged with barium (Ba) and potassium (K) to recover 90.7% of p-DEB at the temperature 450 K. and the pressure 136 atm, wherein toulene was used as the desorbent; Lu and. Lee, in their article, Ind. Eng. Chem., Res. 26, 2024 (1987), tested several adsorbents and desorbents for separation p-DEB from mixed DEB and found that the CDZ zeolite was the most appropriate one to obtain p-DEB with a purity over 95% in the temperature range from 393-463 K. and in the pressure range from 1.5 to 3.0 atm. wherein o-xylene was used as the desorbent.\nSantacesaria, et al. in their article, entitled \"Separation of Xylenes on Y Zeolites in the Vapor Phase. 1. Determination of the Adsorption Equilibrium Parameters and of the Kinetic Regime\", Ind. Eng. Chem. Processes Des. Dev. 24, 78-83 (1985) reported the separation of xylene isomers under gas phase operation was found to be superior to that under liquid phase operation. The improvement was attributed to higher mass transfer rate in vapor. But for both operations a desorbent is required. Some of C.sub.8 aromatics were the most commonly employed desorbent.\nSeveral methods of separation using desorbent as mentioned above require using aromatic compounds as a desorbent, therefore, a distillation operation is thereby needed in order to separate desorbent from product, which is energy intensive and increased separation cost.\nTAN, Chung-Sung, one of the inventors of the present invention, and Tsay, Jeng-Leei, in their article, entitled \"Separation of Xylene Isomers on Silicalite in Supercritical and Gaseous Carbon Dioxide\" Ind. Eng. Chem. Res., Vol. 29, 502-504 (1990) reported an experimental study of the separation of an equal amount of p- and m-xylenes on silicalite using carbon dioxide as the carrier and the desorbent. The results showed that the operations in the gaseous phase carbon dioxide offered a better separation efficiency over those at supercritical conditions. The effects of temperature, pressure and flow rate on the effectiveness of separation were also examined. It was found that, for a pulse of 1.0 cm.sup.3 of xylene isomers and 39.5 g of silicalite, the most appropriate operating conditions were a temperature of around 358 K., a pressure of 476 atm and a flow rate of about 15.0 cm.sup.3 /min.\nTAN, Chun-Sung, in his U.S. Pat. application Ser. No. 07/801,531 filed on Dec. 2, 1991 (which was allowed on Dec. 24, 1992), disclosed a process for separating ethylbenzene and p-xylene from xylene isomers mixture, wherein a compressed high pressure vapor phase CO.sub.2 was used as the carrier. A fixed amount of xylene isomers was fed into a silicalite adsorbent bed for adsorption operation, and a first stage effluent stream containing a rich m-xylene and/or o-xylene, but having substantially no ethylbenzene and p-xylene was obtained from the adsorbent bed, wherein when ethylbenzene or p-xylene contained in the mixture started to appear in the first stage effluent stream, a supercritical CO.sub.2 used as a desorbent was introduced into the adsorbent bed and a second stage effluent stream containing a rich ethylbenzene and p-xylene, but having substantially no other xylene isomers is thereby obtained. Preferably, the first and the second effluent streams could be separately sent into different activated carbon adsorbent beds to adsorb isothermally and isobarically the said xylene isomers products, and the resulting substantially pure CO.sub.2 exiting from said activated carbon adsorbent bed could be recovered and repeatedly used. Alternatively, the introduction of said high pressure gaseous CO.sub.2 could continue until m-xylene, o-xylene and p-xylene were completely separated from the mixture, and then a supercritical CO.sub.2 was introduced as a desorbent such that a second stage effluent stream containing a rich ethylbenzene (EB) was obtained.\nThe object of the present invention is to provide a process for separation of diethylbenzene isomers on silicalite in high pressure gaseous CO.sub.2 or propane to produce a high purity of p-diethylbenzene.\nAnother object of the present invention is to provide a process for separation of diethylbenzene isomers on silicalite in high pressure gaseous CO.sub.2 or propane to to produce a high purity of p-diethylbenzene with a high productivity, that is, the amount of the adsorbent used is relatively lower.\nThe other object of the present invention is to provide a process for separation of diethylbenzene isomers on silicalite in a high pressure gaseous CO.sub.2 or propane to produce a high purity of p-diethylbenzene, wherein CO.sub.2 or propane is used as a carrier in the first stage of the process and as a desorbent in the second stage of the process, and by reducing the pressure of the effluent stream from the adsorbent bed the diethylbenzene isomers products could be separated from the carrier or desorbent."}
{"text": "The present invention relates to the structure of a sawdust discharge port of a chain saw.\nGenerally, a cover member for the body of a chain saw is made relatively thin as a whole for the purpose of causing a decrease in weight of the chain saw. For this reason, it is necessary, with regard to a side wall thereof particularly at the side of sawing section, to reinforce its lower end edge portion constituting a sawdust discharge port and thereby prevent the side wall from being broken or damaged. To this end, the side wall was conventionally reinforced by being made thicker at its lower end portion in such a manner that its outer surface is allowed to protrude outwardly. By so doing, the inner surface of the side wall is made flat at its lower end portion to prevent a turbulency from occurring in the fluid discharged, thereby making the sawdust discharge smooth. This structure, however, had drawbacks in that, during the sawing operation, the above-mentioned outward protrusion of the side wall at the lower end portion thereof hinders the performance of the operations, or sawdust is very likely to be widely scattered into outside the body of the chain saw."}
{"text": "The invention relates to the field of encapsulating nanomaterials and, in particular, to the encapsulation of semiconducting nanowires and carbon-encapsulated semiconducting nanowires."}
{"text": "This invention relates to a method for the combustion of coal in a cyclone combustor, and more particularly, to a method which minimizes the emission of ash particles and other pollutants from such a combustor.\nThe apparatus known as a cyclone coal combustor consists of a cylindrical chamber into which pulverized coal is injected and centrifuged to the cylindrical wall of the chamber by a high velocity, swirling gas flow. Heretofore, most cyclone combustors have been designed for combustion to take place near or at the cylindrical wall, and at a temperature sufficiently high, on the order of about 3000.degree. F., to melt the coal ash. Conventionally, the melted coal ash, or slag, has been drained continuously from the combustor. Although there are two general types of cyclone combustors, the horizontal cyclone of interest in connection with the present application and the reverse flow vertical cyclone, most prior commercial experience is with the horizontal cyclone.\nCommercial horizontal cyclone combustors heretofore used have removed from 70 to 85 percent of coal ash as slag. These commercial cyclone combustors, typically of 100 to 700 million BTU/per hour capacity, were used extensively in large utility and industrial boilers in the 1950's and early 1960's. Their inability however, to control the emission of NO.sub.x, led to the discontinuance of their use in the late 1960's, as environmental concerns grew. Renewed interest arose in the cyclone combustor in the mid 1970's due to the need for a low ash output combustor in conjunction with the U.S. Department of Energy's magnetohydrodynamics (\"MHD\") program. This work, as well as greatly expanded research and development efforts on coal combustion phenomena, led to an improved understanding of the cyclone combustor (C. S. Cook, et al., \"Evaluation of Closed Cycle MHD Power\", DOE Contract Report No. DE-AC01--78ET 10818, Nov., 1981) (\"Ref. 1\").\nA one million BTU/hr. air cooled, horizontal cyclone coal combustor, which was tested as part of an MHD program activity, had a central oil gun located at the center of the closed end of the unit. The oil gun was used to preheat the combustor wall, which was ceramic lined, and to start coal combustion. Pulverized coal, in this combustor, was transported by a primary air stream in about a one-to-one air/coal mass ratio, and injected into the combustor in an annular region surrounding the oil gun. The ceramic liner was maintained at a temperature of about 2200.degree. to 2500.degree. F., a temperature high enough to keep the coal slag in a liquid free flowing state. The ceramic liner was cooled by an air stream which was also used as the secondary air stream. The cooling/secondary air stream was injected at the closed end of the combustor to produce the characteristic swirling gas motion around the axis of the combustor, and the injection velocity was high enough to obtain tangential gas velocities of several hundred feet per second at the cylindrical walls of the combustor. In the above-described apparatus, efforts were made to burn as much of the coal particles as possible in suspension near the cylindrical wall of the combustor, and the unit was operated with conventional pulverized coal, having a conventional particle size distribution. Nevertheless, the above-described prior art device achieved slag retention values in excess of 80 percent and, in some cases, when operated with a somewhat more coarse coal particle size distribution (70% through 200 mesh), slag retention above 90 percent was evident. Significant NO.sub.x emission reductions were obtained with this apparatus and method by operating it at sub-stoichiometric conditions in the range of 60-90% of the stoichiometric air/fuel ratio (Ref. 1). In addition, pulverized limestone was injected as a sorbent directly through the closed end of the combustor, and this resulted in 25 to 35 percent reductions in the SO.sub.2 emissions from the combustor. The reacted sorbent was embedded in the slag and removed with it through a slag tap located at the down stream end of the combustor.\nHorizontal cyclone combustors which were in commercial use in the United States and Germany in the 1950's and '60's had the primary air (which carried the coal) and the secondary air (which produced the swirl) injected tangentially along most of the top axial length of the unit. (Babcock & Wilcox Co., STEAM, ch. 10 (1978 ed.) (\"Ref. 2\")). The U.S. units typically operated with coarse crushed coal (50 percent through 20 to 40 mesh), in order to conserve pulverization power and to provide a method for burning coals having low ash fusion temperatures. Due, however, to the use of such very large coal particles, it was assumed that most of the coal combustion process took place on the surface of the slag, so that excess air, scrubbing the slag layer, was necessary for complete combustion. This, of course, produced high NO.sub.x emissions, and ultimately led to discontinuance of the use of such combustors.\nThe German commercial horizontal cyclone coal combustors were similar in design to the above-described U.S. units. (H Seidl, \"Development & Practice of Cyclone Firing in Germany\", Proc. Jt. Conf. on Combustion ASME-I.Mech.Eng., MIT, Cambridge, Mass., June, 1955, p. 92) (\"Ref. 3\"). However, in the German case, the intended application was the combustion of very high ash (up to 40 percent) coal which had low volatile matter content (under 20 percent). Consequently, a finer, but still relatively coarse pulverized coal particle distribution (approximately 60 percent through 100 mesh, i.e., 150 micron diameter coal particles) was used. These units obtained up to 85 percent slag retention, as compared to only 70 percent in the U.S. units.\nAnother class of cyclone combustors is the vertical, reverse flow type, whose design principles are very similar to those of conventional cyclone dust separators. A detailed experimental study of such a device was performed by Hoy (H. R. Hoy, et al., \"Some Investigations with a Small Cyclone Combustor\" in Jrn. Inst. Fuel, Oct., 1958, p. 429) (\"Ref. 4\"). In the unit studied by Hoy, which had a relatively large 20 million BTU/per hr. coal energy throughout, 80 to 85% of the slag was retained in the combustor. It was observed, however, that the reverse flow had a tendency to re-entrain slag from the walls in some cases.\nOne observation of importance in relation to the present invention should be made here in connection with the above-described cyclone combustors and other cyclone combustors being developed for MHD applications. In this regard, it is of significance that in these units coal was injected very close to the hot, liquid slag-covered walls, and as a result, the coal particles impinged very rapidly on the slag-covered walls. Thus, in such units, even the smaller particles that remained in suspension in the gas stream for but brief periods were in a gas temperature environment in the 3000.degree. F. range. Under these conditions, convective heating of the particles results in rapid pyrolysis and devolatilization of the coal particles. In the cyclone combustors developed for MHD applications, discussed below, in which different coal injection techniques were used, the use of very high temperature air preheat (in the 2000.degree. to 3000.degree. F. range) produces a similar effect. By contrast, in the air-cooled cyclone combustor of the present invention, coal injection occurs in a relatively cold gas environment, and therefore, considerably longer time periods are required and employed for coal pyrolysis.\nThe MHD program gave rise to a need for a high slag retention combustor. Several investigators designed and tested different versions of cyclone combustors intended for the MHD application. The Pittsburgh Energy Technology Center designed and tested a vertical unit similar in concept to Hoy's device. (W. S. Lewellen, et al., \"Modeling Two Phase Flow in a Swirl Combustor\", Aeronautical Research Rpt., (00-4062-5 (1977)) Princenton, N.J.) (\"Ref. 5\"). TRW designed a horizontal unit, similar to the early United States and German units described above. (J. A. Hardgrove, \"MHD Cost Fired Combustor Dev.\", 9th Energy Tech. Conf. Proc., Wash., D.C. Feb., 1982) (\"Ref. 6\"). The TRW unit, however, differed in that it used axial coal injection, as distinguished from the tangential injection used in the commercial units.\nOne significant operational difference between the processes performed by the cyclone combustors intended for the MHD application and the U.S. and German commercial units is the extremely rapid devolatilization that occurs in the MHD units due to the very high temperature air preheat which is used. This condition, coupled with the general use of relatively fine coal particle size distributions in such units appears to have resulted in the MHD units not only in rapid devolatilization, but also very rapid char gasification of much of the fuel while the coal particles were in suspension in the gas stream. These conditions account for very high carbon conversion, but low slag retention (in some cases as low as 30 percent) (Ref. 6).\nAnother prior art technique for the control of certain emissions in the combustion of coal is the injection into combustion chambers and furnaces of limestone or similar calcium oxide compounds as sorbents or binders for sulphur compounds. In the case of the million BTU/hr. combustor described above, 25 to 35 percent reductions in SO.sub.2 emissions were observed with the injection of pulverized limestone, and the reaction products of the limestone were removed with the coal slag. No explanation for this capture process was given (Ref. 1). Generally, the physical states of the sulphur capture process using calcium oxide compounds is by now well understood. The first step in the process is calcination, wherein CaCO.sub.3 is converted to CaO by removal of CO.sub.2 from the CaCO.sub.3 particle. A porous structure is left after calcination. With excess oxygen, sulphur capture leads to the formation of CaSO.sub.4. Under equilibrium conditions, this compound moves toward dissociation above about 2000.degree. F. Sulfation takes place heterogeneously by SO.sub.x contact with CaO, and it is affected by SO.sub.x diffusion through the CaO pore structure. Eventually, however, a layer of CaSO.sub.4 encapsulates the particle and stops the reaction.\nFor cyclone combustor applications, the average gas temperature is 3000.degree. F., which is much too high for equilibrium sulphur capture. However, it is now hypothesized in accordance with this invention that sulphur capture takes place during the time in which the CaCO.sub.3 particle is suspended in the combustion gases, the time for sulphur capture is in the 100 milliseconds time range.\nStill another mechanism by which pollutants are removed from cyclone combustors utilizes the slag layer. The slag removes the mineral matter from the coal in liquid form, and serves as well as a base upon which one can burn up the remaining particles of char in the coal which floats on the slag and to remove the calcium-sulphur compounds resulting from the limestone injection. One problem, however, with prior art techniques has been the reevolution of SO.sub.2 from the slag layer. This occurs because at the temperatures of the slag layer, approximately 2200.degree. to 2500.degree. F., calcium sulphate in the slag will melt and react with species in the slag such as iron compounds or gases above the slag layer such as 0.sub.2 or CO.sub.2. For partial pressures of oxygen of less than 0.1 atmospheres, a condition likely to exist in a cyclone combustor, the rate of SO.sub.2, evolution from slag is in the time range of 15 to 20 minutes. It is, therefore, essential to remove the slag in a time less than this to avoid reevolution of sulphur. Reevolution is retarded somewhat, it has been found, by maintaining local reducing conditions above the slag. However, as noted, compounds such as those of iron can catalyze the reaction which converts chemically bound sulphur in the slag to gaseous form."}
{"text": "(1) Field of the Invention\nThe present invention relates to ash formation during the burning of coal and more particularly to methods and compositions for treatment of coal to reduce the amount of ash deposition onto surfaces during the burning of coal.\n(2) Description of the Related Art\nSub-bituminous coal of the Powder River Basin of the United States typically includes a significant amount of calcium bound within the coal structure. In fact, the typical calcium level of this type of coal burned in industrial boilers in the United States today is substantially higher than it had been in the past and that level is expected to increase in the future as industries continue to turn to lower sulfur level coal.\nWhen the coal is burned, the calcium in the coal is converted to calcium oxide. The formation of calcium oxide results in an ash that is reflective and whiter than the fly ash produced upon combustion of bituminous coal. This reflective ash accumulates on surfaces situated in the structure in which the burning takes place. Such structures will be referred to herein as xe2x80x9cfurnaces,xe2x80x9d as such term is considered in its broad sense to refer to any enclosed structure in which heat is produced. A particular situation in which such ash formation is encountered is in furnaces employed in boiler systems, but the furnaces contemplated herein are not limited to such systems and may be incorporated into any number of uses.\nProminent among the surfaces on which reflective ash tends to accumulate are the furnace tube walls through which heat is to be transferred from the combustion taking place in the furnace. Such ash accumulation is undesirable because the layer it forms over the surfaces is an insulative barrier that reduces the heat transfer through the surfaces, thereby reducing the efficiency of heat transfer from the furnace. Such ash accumulation is also undesirable because the reflective ash layer reflects the heat back into the burner area, increasing the gas and flame temperatures beyond that for which the furnace was designed, which in turn causes the increased heat to radiate back to the fly ash, eventually creating a slagging environment. Moreover, because of the inadequate heat transfer to the water flowing through the furnace wall tubes, the furnace exit gas temperature (FEGT) rises above the design level, increasing the fouling propensity in the convective zone and, in the case of the boiler, finally increasing the boiler exit gas temperature. The increased FEGT also raises the temperature of steam in downstream heat absorption sections above design conditions, requiring use of cooling spray water to reduce the steam temperature.\nThe formation of this type of ash has become more pronounced in recent years. Many boiler furnaces were designed for burning high sulfur bituminous coal. However, as alluded to above, beginning in the late 1970\"\"s and early 1980\"\"s, environmental concerns led to conversion from burning high sulfur bituminous coal to burning low sulfur coals, such as that from the Powder River Basin in Wyoming (PRB coal), began. Even though the ash content of PRB coals is lower than that of the high sulfur coals they replace, PRB coals tend to be high in calcium. Thus, burning the lower sulfur coals in the furnaces designed for relatively high sulfur coal has resulted in increased slagging, and particularly increased white ash formation. See, for example, xe2x80x9cPRB Coal Switch Not a Complete Panacea,xe2x80x9d by Buecker, B. and Meinders, J., Power Engineering, November 2000, pp. 76-80.\nConventionally, equipment such as soot blowers and water lances have been employed to reduce slagging and lower the FEGT, but with limited success and the additional costs, efforts and interference associated with such equipment. Moreover, use of a water lance is undesirable because it introduces cold water into the furnace, inducing thermal stress to the tubes, decreasing the furnace wall tube life and increasing the maintenance and replacement costs of the boiler.\nOther prior art methods have addressed the problem of ash accumulation on furnace walls with chemical techniques for darkening the ash on the walls. For example, U.S. Pat. No. 5,819,672, incorporated herein by reference, describes the addition of a darkening agent such as iron oxide and, preferably, also a fluxing agent to produce a dark ash coating or an additional dark ash coat over the existing ash on the furnace walls. Because the ash is darkened, not only is the tendency of the ash to reflect heat back into the furnace reduced, but the heat absorption by the ash is increased, thereby reportedly aiding transfer of heat from the furnace through the walls thereof. Although the additives may be applied to the coal, the preferred method contemplates applying the dark coat directly to the ash. In any event, such techniques do not eliminatexe2x80x94or even reducexe2x80x94ash accumulation and suffer from various other disadvantages as well. For example, pursuant to such techniques, ash still is allowed to build up on the surfaces at previous rates with the attendant problems, such as the need for routine cleaning or replacement.\nIn the above-noted article in Power Engineering, it is reported that ADA Environmental Solutions has experimented with the application of a proprietary mixture of iron oxides and stabilizing chemicals to coal prior to combustion to enhance the viscosity characteristics of the slag formed from burning PRB coal and that the preliminary results from this experimentation xe2x80x9chave been very promising.xe2x80x9d However, no further information is provided in the article as to the composition of the additive and, although the article later discusses ash control, the article nowhere discusses the additive with respect to the ash control. Indeed, the article indicates that ash accumulation and high FEGT are still significant problems, requiring the use of water lances.\nThus, the industry is still searching for an effective and efficient means for darkening the ash and reducing the FEGT, slagging and ash accumulation on surfaces in coal-burning furnaces. Techniques that accomplish such objectives, while avoiding the need for purchasing and operating equipment such as soot blowers and water lances would be particularly desirable. And, of course, it is also especially desirable that the technique avoid raising adverse environmental implications.\nBriefly, therefore, the present invention is directed to a novel method for inhibiting accumulation of reflective ash on surfaces in a furnace in which high calcium-containing coal is burned. According to the method, an effective amount of an iron compound is added to the coal to produce treated coal, free of added fluxing agent, and the treated coal is then burned.\nThe present invention also is directed to a novel method for increasing the melting point of ash produced during the burning of calcium-containing coal. According to the method, an effective amount of an iron compound is added to the coal to produce treated coal, and the treated coal is burned, producing ash of increased melting point.\nAmong the several advantages found to be achieved by the present invention, therefore, may be noted the provision of a method for darkening ash formed in the combustion of coal; the provision of such method that also reduces the tendency of the ash to accumulate on surfaces in the furnace; the provision of such method that also reduces the FEGT in the furnace; the provision of such method that improves the overall boiler efficiency and reduces generation cost; the provision of such method that eliminates the need for soot blowers and water lances; the provision of such method that reduces slagging in the furnace; and the provision of such method that avoids introduction of adverse environmental consequences.\nIn accordance with the present invention, it has been discovered that, surprisingly, a darker ash may be produced and accumulation of reflective ash on surfaces in a furnace may be reduced simply by adding iron oxide to calcium-containing coal prior to combustion of the coal in the furnace. Thus, rather than lowering the melting point of the ash or calcium components thereof by addition of a fluxing agent to encourage adhesion to the furnace surfaces as described in U.S. Pat. No. 5,819,672, the process of the present invention involves raising the melting point to inhibit such adhesion. Moreover, it has been found that the addition of the iron oxide not only darkens the ash and inhibits the tendency of the ash to adhere to the furnace surfaces but, accordingly, also reduces the FEGT and consequently slag and fouling deposit formation. In fact, the improvements in furnace performance resulting from the method of this invention have been discovered to be even greater than those achieved with the conventional treatments by soot blowers and water lances, thus eliminating the need for such equipment. And all of this has been found to be accomplished without any detected adverse environmental consequence.\nWhile not wishing to be bound by any particular theory of operation, applicant believes that the present invention operates by conversion of the calcium oxide formed upon combustion of the coal to calcium ferrite, thereby converting calcium oxide in the ash to relatively higher melting point and darker calcium ferrite. This conversion may be illustrated by the following idealized formula wherein ferric oxide is the additive:\nCaO+Fe2O3xe2x86x92Ca(FeO2)2 \nIn any event, according to the present invention, a composition comprising an iron compound is added to the calcium-containing coal prior to combustion of the coal, preferably prior to delivery of the coal to the furnace and most desirably prior even to grinding the coal. The iron compound in the most desirable embodiment contemplated by the invention is iron oxide, especially ferric oxide. It is also desirable that the additive composition contain no other component that interferes with the ability of the iron compound to raise the melting (or fusion) point of the resulting ash. The additive composition particularly should not contain a fluxing compound or an adhesive or other substance that increases the tendency of the ash to adhere to the furnace surfaces.\nIn one preferred embodiment, the additive is a composition of iron oxide (especially in the form of ferric oxide) and clay, which is primarily silica with traces of alumina and other calcium and magnesium compounds. One such preferred formulation comprises 93% by weight ferric oxide, 5% by weight silica, with the remainder made up of alumina and other calcium and magnesium compounds. Hematite ore has been found to be a particularly appropriate additive composition. However, in another preferred embodiment, the additive composition may be the iron compound, with no other ingredients other than at most minor impurities.\nThe preferred form of the additive composition is a powder. However, other forms, such as a suspension of that powder in a liquid such as a liquid hydrocarbon (e.g., kerosene), may be employed if so desired. Although the liquid may be water, such is not desirable for the obvious thermodynamic and other disadvantages of introducing water into the combustion process.\nThe additive composition is applied to the coal, such as by spraying or spreading, in an amount sufficient to provide an effective amount of the iron compound to combine with the calcium content of the coal. In the context of inhibition of ash on the furnace walls, by xe2x80x9can effective amountxe2x80x9d of the iron compound what is meant is an amount that is sufficient to result in less ash deposition on the furnace walls when the coal treated with the iron compound is burned than forms on the walls when equivalent coal without the iron compound treatment is burned under equivalent conditions. In the context of increasing the melting point of the ash produced when the coal is burned, by xe2x80x9can effective amountxe2x80x9d of the iron compound what is meant is an amount that is sufficient to increase the melting point of the ash produced when the coal is burned over the melting point of the ash produced when equivalent coal without the iron compound treatment is burned under equivalent conditions. Such ash having a melting point higher than that of the ash produced when equivalent coal without the iron compound treatment is burned under equivalent conditions is referred to herein as xe2x80x9cash of increased melting point.xe2x80x9d\nThe optimal amount of iron compound to be added to the coal depends on the calcium content of the coal. Generally, however, when the iron compound is ferric oxide and the coal is of sub-bituminous type from the Western United States and particularly PRB coal, the optimal amount of iron compound has been found to be from about 0.1% to about 1.0%, such as about 0.1% to about 0.75%, more preferably about 0.25% to about 1.0%, even more preferably about 0.25% to about 0.75%, especially about 0.5%, based on the weight of the coal. Based on the theorized formula set forth above, this represents, surprisingly, only about one-sixth the amount of ferric oxide required by the stoichiometry. Although greater amounts of iron compound may be used, it is believed that there is currently no economic advantage to doing so.\nAfter application of the additive composition to the coal, the treated coal then is ground, if not already ground, and conveyed to the furnace, wherein it is exposed directly to the flame envelope of the furnace combustion process. As described above, the resulting ash is darker and has a higher melting point compared to the ash formed from coal not treated in accordance with this invention. In fact, rather than tenaciously adhering to the exposed surfaces, the darkened ash is gas-borne fly ash, most of which escapes from the furnace with the exhaust gases, reducing or even eliminating the need for cleaning out the ash from the furnace wall. And because the product comprises calcium, iron and oxygen, which pose no environmental concern. Moreover, because the FEGT is lower when the coal is treated with the iron compound according to this invention, use of the water lance may be eliminated.\nThe following examples describe preferred embodiments of the invention. Other embodiments with the scope of the claims herein will be apparent to one skilled in the art from consideration of the specification or practice of the invention as disclosed herein. It is intended that the specification, together with the examples, be considered exemplary only, with the scope and spirit of the invention being indicated by the claims which follow the examples."}
{"text": "1. Field\nThe present invention relates generally to a manufacturing method of an organic light emitting diode (OLED) display.\n2. Description of Related Art\nDisplay devices display images, and recently, an organic light emitting diode display has been in the spotlight.\nThe organic light emitting diode display has a self-emitting characteristic and does not need a separate light source such that the thickness and weight thereof are reduced in comparison to a liquid crystal display. In addition, the organic light emitting diode display has low power consumption, high luminance, high response speed, and the like.\nIn general, the OLED display includes a substrate, a plurality of organic light emitting diodes disposed on the substrate for displaying an image, and an encapsulation member facing the substrate while encapsulating the organic light emitting diodes therebetween.\nRecently, an OLED display using a thin film encapsulation layer as an encapsulation member has been developed.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the present invention, and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "1. Field\nThe present application relates to a fuel-level measuring apparatus that measures the fuel level in a fuel tank installed in a vehicle.\n2. Description of the Related Art\nIn a vehicle that uses gasoline or light oil as fuel, it is necessary to indicate the fuel level in a fuel tank so that the driver can know when the tank needs to be refilled and how far the car can drive.\nA known fuel-level measuring apparatus that measures the fuel level in a fuel tank is disclosed in, for example, Japanese Unexamined Patent Application Publication No. 2013-50410. The apparatus disclosed therein includes a float provided in a fuel tank, the float moving up and down with the surface level of the fuel, and, by making a float arm of the float slide over a resistor plate and converting the surface level into an electric potential difference, the apparatus measures the height of the surface and indicates the fuel level calculated from the height of the surface to a driver.\nHowever, the technique disclosed in Japanese Unexamined Patent Application Publication No. 2013-50410 has the following problems: because the fuel-level measuring apparatus, such as the float, is disposed in the fuel tank, the capacity of the fuel tank decreases by an amount corresponding to the volume of the fuel-level measuring apparatus; because the fuel-level measuring apparatus needs to be disposed in the fuel tank, the shape of the fuel tank is restricted and cannot be thin or complicated; and, because the surface level of the fuel in the tank changes significantly depending on the orientation and driving conditions of the vehicle, it is impossible to measure the accurate fuel level."}
{"text": "Users may interact with various types of media objects comprising a wide range of information. In an example, a user may view a beach image shared by another user through an image sharing service. In another example, the user may read a videogame console article published through a videogame website. In another example, the user may watch a live video stream of a soccer game through a video streaming app. In another example, the user may listen to a podcast through an online radio station. Many content providers may spend substantial amounts of time and computing resources attempting to combine content (e.g., a message, a link, an image, an interactive game, etc.) and media objects (e.g., an image, video, and/or audio). For example, a website may play an initial video before playing a requested video (e.g., a videogame website may force users to watch a video about a brief history of the website before being able to watch a videogame review show). In another example, a news website may display an image overlaying a news article. Unfortunately, the combination of content and media objects may be arbitrary, random (e.g., subject matter of the image may have little to do with subject matter of the article, such as a weather image overlaying an article about a corrupt politician), and/or annoying (e.g., the image may obscure part of the article). Thus, manual and human intervention are required to improve the combination of content and media objects so that they better match."}
{"text": "This invention relates to a novel integrated-circuit RC oscillator. More specifically, this invention relates to an RC oscillator which is entirely fabricated on an integrated-circuit chip, and has an output oscillation frequency which may be externally trimmed.\nOscillators are used in microprocessors to convert energy from a DC source into a cyclic energy which is usually supplied to the microprocessor clock. Crystal oscillators are sometimes employed to produce the desired oscillation frequency. While crystal oscillators produce good quality oscillating voltage signals, they are more expensive than most oscillators. Therefore, other types of less expensive oscillators have been used. In the past, RC oscillators were made having the passive resistive and capacitive elements positioned off the integrated-circuit chip because the elements occupied a considerable amount of semiconductor area. As semiconductor technology has advanced, it has been possible to position some of the resistive or capacitive elements on the IC chip.\nIt is oftentimes desirable for a user to be able to change the output voltage signal frequency of an oscillator after purchasing the oscillator. In order to do so, a user would purchase a new IC chip having an oscillator which produced the desired oscillation frequency. However, this method of modulation is expensive and inefficient. RC oscillators were then developed which included some resistive and capacitive elements off-chip so that the user could adjust the oscillating frequency. These oscillators were not entirely satisfactory, however, because the off-chip elements utilized too much space.\nThe difficulties in the preceding are not intended to be exhaustive but rather are many which may tend to reduce the usefulness of prior art RC oscillators. Other noteworthy problems may also exist; however, those presented above should be sufficient to demonstrate that RC oscillators appearing in the past will admit to worthwhile improvement.\nIt is therefore a general object of the invention to provide a novel RC oscillator which will obviate or minimize difficulties of the type previously described.\nIt is a specific object of the invention to provide an RC oscillator which produces an output voltage signal at an oscillating frequency which may be easily adjusted.\nIt is another object of the invention to provide an RC oscillator which minimizes the number of components which are external to the semiconductor chip on which the oscillator is fabricated.\nIt is still another object of the invention to provide an RC oscillator which is inexpensive and which provides a good quality voltage signal which consistently oscillates at a desired frequency."}
{"text": "A liquid crystal display apparatus has been widely used as a display apparatus such as a television or a monitor.\nParticularly, an active matrix liquid crystal display apparatus in which a transistor which serves as a switching element is provided in each of pixels has been recently mainstream of liquid crystal display apparatuses, from viewpoints of high display quality and high-speed response.\nIn general, a liquid crystal display panel constituting the active matrix liquid crystal display apparatus is made by combining, as a pair, a color filter substrate serving as a counter substrate with an active matrix substrate.\nThe active matrix substrate includes pixel electrodes arranged in a grid pattern, a plurality of gate lines and a plurality of source lines, and thin-film transistors (TFT) each serving as a switching element and each provided in the vicinity of an intersection of a corresponding one of the plurality of gate lines and a corresponding one of the plurality of source lines. Each of the plurality of gate lines is provided between corresponding two adjacent ones of the pixel electrodes, and each of the plurality of source lines is provided between corresponding two adjacent ones of the pixel electrodes. The plurality of gate lines and the plurality of source lines are provided so as to orthogonally intersect with each other.\nIn a pixel region which is a display region of the liquid crystal display panel, in general, an insulating layer is provided between (i) a gate line and (ii) respective of a pixel electrode and a source line.\nThe plurality of gate lines and the plurality of source lines are provided so as to be partially exposed, for their electrical connections with other components, in a terminal region which is a peripheral region of the liquid crystal display panel.\n(Patent Literature 1)\nA liquid crystal display apparatus having such an arrangement is exemplified by a liquid crystal display apparatus disclosed in Patent Literature 1. Specifically, Patent Literature 1 discloses an arrangement as below. In a display region (pixel region), a gate electrode and a source electrode are formed on a first transparent substrate via a gate insulating film. Further, the gate electrode, the gate insulating film, and the source electrode are covered by a protection film.\nIn a terminal region on the periphery of the display region, the protection film is formed in the terminal region except a region where a line terminal section is provided.\n(Patent Literature 2)\nPatent Literature 2 discloses an arrangement as below. In a pixel region, an upper insulating film is provided on a signal line (source line) so that a surface of the upper insulating film is planarized. Then, a pixel electrode is formed on the upper insulating film.\n(Patent Literature 3)\nPatent Literature 3 discloses an arrangement as below. In a pixel section (pixel region) and a terminal section (terminal region), an underlayer is provided on a substrate, and then, lines are formed on the underlayer."}
{"text": "The present invention generally relates to an image forming system, and more particularly to a toner supplying device of an image forming system, which device includes a cleaner portion for removing a residual toner from a photosensitive medium and an integrally formed toner supplying portion for supplying toner to a developing unit so that a toner image is formed on a copy sheet.\nIn image forming systems such as facsimile machines, printers and copiers, a cleaner toner magazine (CTM) unit is used in order to remove a residual toner on a photosensitive medium and store the residual toner, and to supply new toner to a developing unit so that a toner image is formed on a copy sheet. FIG. 12 shows a CTM unit used in an image forming system. In FIG. 12, the CTM unit includes a cleaner part 3 for removing the residual toner from a photosensitive medium (not shown) and a toner supplying part 5 for supplying new toner to a developing unit (not shown); these two parts are built into the one-piece unit. The cleaner part 3 includes a cleaning blade 1 for removing the residual toner, and a compartment 2 for storing the residual toner which is received from the photosensitive medium due to the cleaning blade 1. The toner supplying part 5 includes a toner transport screw 4 for supplying new toner to the developing unit.\nIn an image forming system of a different type having no CTM unit, the photosensitive medium, the developing unit, a replenishment toner storage tank, and a used toner storage tank are permanently united into a one-piece unit. When the replenishment toner in the storage tank is used up, it is necessary to replace the photosensitive medium or the developing unit, which is still usable and has not reached the end of its operating life. The image forming system with the CTM unit as described above has an advantageous feature in that the CTM unit, singly, can be replaced with a new CTM unit when the toner is used up, thus allowing the running cost per copy to be lowered.\nHowever, when the CTM unit is replaced, for the replenishment of toner, there is a problem in that the replenishment toner partially sticks to an inside wall of the toner supplying part 5. The sticking of the toner may occur due to the changes of the chemical characteristics of the toner (e.g., the flowability or the adsorption characteristics) under certain environmental conditions or due to the physical arrangement of the CTM unit within the image forming system. For this reason, the amount of usable toner in the CTM unit is unsuitably reduced. Thus, it is desirable to make effective use of the replenishment toner in the CTM unit."}
{"text": "The present invention relates to a highly reliable male/female fluid coupling with an improved metallic seal for sealably maintaining fluid pressure after repeated mating and unmating operations and, more particularly, relates to an improved metallic seal coupling which allows for both axial and radial deviation between the coupling halves during the mating operation.\nConventional male/female couplings have long been used to sealingly join pressurized fluid lines. The coupling halves, when joined, automatically enable the fluid and/or fluid pressure to be transmitted through the coupling without a significant pressure differential. When unmated, the value stems in the body of each coupler half automatically seal the fluid within the coupler and its respective fluid line. Such couplers preferably are able to reliably maintain high fluid pressure after repeated mating and unmating operations.\nThe seal between the joined coupling halves may conventionally be provided by either an elastomeric or metallic member. The low cost, chemical inert quality and noncorrosive characteristics of elastomeric sealing materials has resulted in the commonplace use of such sealing members for many applications. Such elastomeric seals typically may be manufactured from natural or synthetic rubber, Teflon, nylon, PEEK, or similar materials.\nIn other applications, especially those wherein high reliability is paramount or the contained fluid is highly abrasive, metallic seals are preferable. One such application involves couplers commonly used in subsea petroleum recovery operations, wherein a plurality of coupler halves are mounted on respective control module mounting bases, so that clamping and unclamping of the bases mates and unmates the couplers and thereby interconnects numerous fluid pressure transmission lines.\nOne type of coupler suitable for reliably interconnecting fluid lines employs a metallic face seal. Sealing engagement is provided by a metallic seal sandwiched between surfaces of the coupler halves, with each sealing surface being generally perpendicular to the axis of the coupler. The metallic seal is permanently deformed into sealing engagement with each sealing surface by an axially directed mating force, and thus must be replaced between successive mating operations. A coupling with a reusable metallic radial seal is disclosed in U.S. patent application Ser. No. 908,593, hereby incorporated by reference.\nOne disadvantage of prior art metallic seal couplers is the expense and inherent decreased reliability associated with maintaining close tolerance between the sealing surfaces of the coupler halves. Close axial tolerances must be maintained between the surfaces for engagement with the face seal for the first coupler described above, and similarly close alignment or radial tolerance of the cylindrical surfaces for engagement with the radial seal is necessary for the second coupling described in U.S. patent application Ser. No. 908,593. These tolerance problems, although not particularly troublesome for many coupling operations, are particularly difficult and expensive to correct for remote coupling operations, such as petroleum subsea operations.\nThe disadvantages of the prior art are overcome by the present invention, and an improved coupler is hereinafter described with metallic seals positioned to allow for both axial and radial deviation between the coupling halves."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing an optical waveguide device for widespread use in optical communications, optical information processing and other general optics.\n2. Description of the Related Art\nIn general, an optical waveguide for an optical waveguide device is constructed in such a manner that cores serving as a passageway for light are formed in a predetermined pattern on a surface of an under cladding layer, and that an over cladding layer is formed so as to cover the cores. Such an optical waveguide is typically formed on a surface of a substrate such as a metal substrate and the like, and is manufactured together with the substrate to provide an optical waveguide device.\nA conventional method of manufacturing such an optical waveguide device is as follows. First, as shown in FIG. 6A, an under cladding layer 20 is formed on a surface of a substrate 10. Then, as shown in FIG. 6B, a photosensitive resin for the formation of cores is applied to a surface of the under cladding layer 20 to form a photosensitive resin layer 30A. Next, irradiation light L is directed through a photomask M formed with an opening pattern corresponding to the pattern of cores 30 (with reference to FIG. 6C) toward the photosensitive resin layer 30A. The irradiation light L is caused to reach the photosensitive resin layer 30A through openings of the opening pattern, thereby exposing portions of the photosensitive resin layer 30A thereto. The irradiation light L is directed to the photosensitive resin layer 30A at right angles thereto. A photoreaction proceeds in the portions of the photosensitive resin layer 30A exposed to the irradiation light L so that the exposed portions are hardened. Then, development is performed using a developing solution to dissolve away unexposed portions of the photosensitive resin layer 30A, as shown in FIG. 6C. The remaining exposed portions become the cores 30 formed in a predetermined pattern. The cores 30 are typically rectangular in sectional configuration. Thereafter, as shown in FIG. 6D, a photosensitive resin for the formation of an over cladding layer is applied to the surface of the under cladding layer 20 so as to cover the cores 30, to thereby form a photosensitive resin layer 40A. Next, irradiation light is directed toward the photosensitive resin layer 40A to expose the photosensitive resin layer 40A thereto. This causes the photosensitive resin layer 40A to be formed into an over cladding layer 40. In this manner, an optical waveguide W0 is formed on the surface of the substrate 10, for example, as disclosed in Japanese Published Patent Application No. 2004-341454.\nIn such a conventional method, however, the cores 30 have side surfaces 30a formed as roughened surfaces in some cases, as shown in FIGS. 7A and 7B. Such cores 30 increase the propagation losses of light propagating inside the cores 30 because of leakage of light from the side surfaces 30a formed as the roughened surfaces or other reasons. There is apprehension that the above-mentioned optical waveguide does not function as optical interconnection. When the surface roughening of the side surfaces 30a of the cores 30 is even slight but is greater than the wavelength of the light propagating inside the cores 30, the surface roughening causes scattering losses. Also in this case, the propagation losses increase, and there is apprehension that the above-mentioned optical waveguide does not function as optical interconnection, as described above. FIG. 7B is a view drawn based on a photograph in perspective of a core 30 enclosed with the circle C of FIG. 7A which is magnified 700 times with an electron microscope. By magnifying the core 30 700 times with the electron microscope in this manner, it can be seen that the side surfaces 30a of the cores 30 are formed as the roughened surfaces.\nThe present inventors have made studies to diagnose the cause of the formation of the side surfaces 30a of the cores 30 as the roughened surfaces. In the course of the studies, the present inventors have found that the surface roughening of the side surfaces 30a of the cores 30 occurs either in the case (1) where a metal substrate 11 (with reference to FIG. 8) in common use which is made of stainless steel (SUS) foil and the like is used as the substrate 10 (with reference to FIGS. 6A to 6D) or in the case (2) where a substrate (PET substrate) 12 (with reference to FIG. 9) made of a material having polyethylene terephthalate (PET) as a main component is used as the substrate 10.\nAs a result of further studies on the case (1) where the metal substrate 11 made of stainless steel foil and the like is used, the present inventors have found that the metal substrate 11 made of stainless steel foil and the like includes a roughened surface having an arithmetic mean roughness (Ra) of not less than 0.095 μm. For this reason, in the above-mentioned step of forming the cores 30, the irradiation light L for use in the exposure which passes through the openings of the photomask M is transmitted through the photosensitive resin layer 30A for the core formation and the under cladding layer 20, and thereafter is reflected diffusely from the roughened surface of the metal substrate 11 because of the surface roughening, as indicated by arrows shown in FIG. 8. The diffusely reflected irradiation light L is transmitted through the under cladding layer 20 obliquely upwardly from below. Then, boundary surfaces (surfaces that are to become the side surfaces 30a) for the patterning of the cores 30 are exposed to the diffusely reflected irradiation light L directed obliquely from below, as indicated by the arrows, in future core regions S included in the photosensitive resin layer 30A for the core formation. This exposure to the light directed obliquely from below results from the above-mentioned diffuse reflection, and is uneven. Thus, it has been found that an unwanted photoreaction proceeds unevenly at the surfaces that are to become the side surfaces 30a of the cores 30 because of the exposure to the light directed obliquely from below to result in the increased width of the cores 30 and the formation of the side surfaces 30a of the cores 30 as the roughened surfaces. In other words, the surfaces that are to become the side surfaces 30a of the cores 30 have both portions subjected to a low degree of exposure to light and portions subjected to a high degree of exposure to light because of the diffuse reflection of the irradiation light L. In the subsequent step of development, the portions subjected to a low degree of exposure to light are dissolved away from the surfaces that are to become the side surfaces 30a of the cores 30, and the portions subjected to a high degree of exposure to light remain unremoved. Thus, the side surfaces 30a of the cores 30 are formed as the roughened surfaces.\nThe present inventors have made further studies on the case (2) where the PET substrate 12 is used. As a result of the further studies, it has been found that the PET substrate 12 contains an additive component such as a lubricant material 12a and the like, as shown in FIG. 9, and that the irradiation light L for use in the exposure in the above-mentioned step of forming the cores 30 impinges upon and is reflected from the additive component such as the lubricant material 12a, to thereby follow irregular paths. It has also been found that most of the irradiation light L reaches the bottom surface (the back surface) of the PET substrate 12. That is, in the above-mentioned step of forming the cores 30, the irradiation light L for use in the exposure which passes through the openings of the photomask M is transmitted through the photosensitive resin layer 30A for the core formation and the under cladding layer 20, and thereafter enters the interior of the PET substrate 12, as shown in FIG. 9. In the PET substrate 12, the irradiation light L impinges upon the lubricant material 12a and the like to follow irregular paths, and then reaches the bottom surface of the PET substrate 12 at an angle. In general, the back surface of the PET substrate 12 is in contact with a mounting surface of a mounting table and the like (not shown) for placing the PET substrate 12 thereon, the mounting surface being impervious to the irradiation light L. For this reason, the irradiation light L reaching the bottom surface of the PET substrate 12 does not exit from the back surface of the PET substrate 12 but is reflected from the bottom surface of the PET substrate 12 obliquely upwardly as indicated by arrows in FIG. 9. Thereafter, the reflected irradiation light L impinges upon the lubricant material 12a and the like in the PET substrate 12 to follow further irregular paths, and then exits from the front surface of the PET substrate 12 obliquely upwardly. The irradiation light L exiting obliquely upwardly is transmitted through the under cladding layer 20 obliquely upwardly from below. Then, the boundary surfaces (the surfaces that are to become the side surfaces 30a) for the patterning of the cores 30 are exposed to the irradiation light L directed obliquely from below in the future core regions S included in the photosensitive resin layer 30A for the core formation. This exposure to the light directed obliquely from below comes from the diffuse reflection resulting from the lubricant material 12a and the like contained in the PET substrate 12 as mentioned earlier, and is uneven. Thus, as in the case (1) where the metal substrate 11 made of stainless steel foil and the like is used, an unwanted photoreaction proceeds unevenly at the surfaces that are to become the side surfaces 30a of the cores 30 because of the exposure to the light directed obliquely from below to result in the formation of the side surfaces 30a of the cores 30 as the roughened surfaces.\nThe assignee of the present application has proposed a method of manufacturing an optical waveguide device which is capable of suppressing the surface roughening of core side surfaces of an optical waveguide when the optical waveguide is formed on a surface of a substrate, and has already applied for patents (Japanese Patent Applications Nos. 2008-130727, 2008-234735, 2008-239346, and 2008-269688).\nJapanese Patent Application No. 2008-130727 discloses a technique for forming an under cladding layer containing an irradiation light absorbing agent on a roughened surface of a metal substrate or forming an irradiation light absorbing layer on the roughened surface of the metal substrate prior to the formation of an under cladding layer. In this technique, irradiation light for the formation of cores is absorbed or attenuated by the irradiation light absorbing agent contained in the under cladding layer or by the irradiation light absorbing layer before or after being reflected diffusely from the surface of the metal substrate. This significantly reduces the amount of irradiation light directed obliquely from below to roughen the surfaces that are to become the side surfaces of the cores by exposing the surfaces thereto. As a result, this technique suppresses the surface roughening of the side surfaces of the cores.\nJapanese Patent Application No. 2008-234735 discloses a technique in which a material that absorbs irradiation light for the formation of cores is used for a substrate having a roughened surface. In this technique, the irradiation light for the formation of the cores is absorbed by the substrate when the irradiation light reaches the surface of the substrate. This significantly reduces the amount of irradiation light directed obliquely from below to roughen the surfaces that are to become the side surfaces of the cores by exposing the surfaces thereto. As a result, this technique suppresses the surface roughening of the side surfaces of the cores.\nJapanese Patent Application No. 2008-239346 discloses a technique in which either a colored-layer-coated substrate having a back surface on which a colored layer of a color that absorbs irradiation light is formed or a colored substrate entirely colored in a color that absorbs the irradiation light is used as a PET substrate. In this technique, the irradiation light for the formation of cores is absorbed by the colored layer or the colored substrate. This significantly reduces the amount of irradiation light directed obliquely from below to roughen the surfaces that are to become the side surfaces of the cores by exposing the surfaces thereto. As a result, this technique suppresses the surface roughening of the side surfaces of the cores.\nJapanese Patent Application No. 2008-269688 discloses a technique in which either only a front surface or both front and back surfaces of a substrate (a silicon wafer, a glass substrate and the like) have an arithmetic mean roughness (Ra) in the range of 1 to 2 nm. In this technique, the surface that reflects irradiation light for the formation of cores is so smooth that the diffuse reflection of the irradiation light from the surface is suppressed. This significantly reduces the amount of irradiation light directed obliquely from below to roughen the surfaces that are to become the side surfaces of the cores by exposing the surfaces thereto. As a result, this technique suppresses the surface roughening of the side surfaces of the cores.\nIn this manner, the above-mentioned techniques of the Applications already filed by the assignee of the present application can suppress the surface roughening of the side surfaces of the cores to reduce the propagation losses of light. However, even when the surface roughening of the side surfaces of the cores is suppressed, the side surfaces of the cores are formed as roughened surfaces in some cases. The present inventors have made studies to diagnose the cause of such a phenomenon. As a result, the present inventors have found that the cause of the phenomenon lies in the contamination of the photomask for use in the formation of the cores. Specifically, a material for the formation of the cores and the like might be deposited onto portions of the photomask which surround the openings of the photomask because of the use of the photomask. In that case, the deposits on the portions of the photomask which surround the openings cause the irradiation light passing through the openings to follow irregular paths. This makes the irradiation light reaching the side surfaces of the cores uneven to result in the roughened side surfaces of the cores.\nIn the formation of an optical waveguide on the surface of a substrate, it has hence been required to suppress the surface roughening of the side surfaces of the cores, thereby reducing the propagation losses of light, regardless of the type of substrate."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for connecting and locking structural panels, and more particularly to a locking structure of floor panels and a method of locking floor panels using the locking structure.\n2. Background Description\nLocking mechanism for floor panels are known, for example, DE 100 08 108. In such floor panels, a groove forms an upper lip and a lower lip over the entire length of at least one lateral border, and a tongue is formed on the opposite lateral border which corresponds to the groove. To lock panels that have been connected, the tongue is provided with at least one projection and the groove exhibits at least one indentation on the lip that faces the projection. The locking mechanism created by the corresponding profiling of the tongue and the groove forms a so-called “click profile”. Two panels are connected by inserting the tongue of one panel into the groove of the other panel, until the projection on the lower side of the tongue clicks into the recess formed in the lower lip of the groove.\nIn addition, panels are known from U.S. Pat. No. 4,426,820 or CH 562 377 A5 which are connected in the laying process by inserting the tongue of a new panel in an angular fashion into the groove of an already laid panel, and then lowering the panel to the floor. The tongue and groove of the two panels engage with each other so that panels are locked both in the transverse direction and perpendicular to the upper panel surface. These types of floor panels, as well as others, can be manufactured from diverse materials such as, for example, from a wood material, particularly MDF or HDF. The upper side is formed with a decorative layer applied to the core.\nPanels that are combined to form a floor cover must be tight at the area of connection in the top surface to prevent moisture from penetrating the connection and to prevent the core material of the panels from swelling. For this reason the “clicking” mechanisms are designed in such a way that two panels inserted into each other are in a state of tension and the connected points are solidly pressed together at the surface. However, the core material of known panels is relatively sensitive, which is often not realized, particularly by those engaged in home repair. If care is not taken when laying the panels and if strong forces are exerted when a new panel is inserted or swung downwards there is the danger that the lower lip of the groove will buckle. This results in the required stress not being able to be applied to the floor panels during installation since the position of the interlocking mechanism has changed. Accordingly, it will become impossible to firmly join the panels. Also, if the lower lip buckles too much it is possible that the groove will be broken, a circumstance that cannot be identified from the outside.\nIt is further noted that the production of the click profiles is rather involved, since maintaining the given tolerances requires careful milling of the core material and demands continuous quality control. If, for example, tool wear results in the projection on the tongue being insufficiently high, the two panels will not lock together in solid fashion, and the danger will arise of the panels separating or of surface gaps forming at the connecting points. If the projection is too high it will not fit the lower lip of the panel provided with the groove. The danger may also arise of the lower lip buckling and being broken or the projection being damaged (ground down), making a solid connection of the panels impossible."}
{"text": "Abbreviations used herein shall have the following meanings:\nAF Amplify-and-Forward\nCRC Cyclic Redundancy Check\nDCI Downlink Control Information\nDF Decode-and-Forward\nDFTS Decode and Forward TS\nDL Downlink\nFSR Frequency Selective Relay (Repeater)\nFTR Frequency Translating Relay (Repeater)\nOFDM Orthogonal Frequency Division Multiplex\nOFR On-Frequency Relay (Repeater)\nPDCCH Physical Downlink Control CHannel\nCRNTI Cell Radio Network Temporary Identifier\nTTI Transmission Time Interval\nUL Uplink\nOne recent development within modern telecommunication standards is the so-called Long Term Evolution (LTE) radio interface, and a further development of the LTE, namely LTE Advanced. These are both OFDM based systems.\nOne of the most important improvement areas in the so-called LTE-Advanced technology is the increase of data rates available for users at the cell edge and indoor. A very promising technique to achieve high data rates in such difficult locations is the deployment of relays. Relays are usually classified into Layer 1 (L1), Layer 2 (L2), and Layer 3 (L3) relays depending on which Open System Interconnection (OSI) layer they operate on. The OSI model is a conceptual model for telecommunication consisting of 7 layers (physical, data-link, network, transport, session, presentation, and application) Note however that the different layers refer to the user plane of the relay node and a L1 relay may use e.g. L3 control plane signaling.\nL1 relays are commonly denoted Amplify-and-Forward (AF) relays or sometimes equivalently repeaters. An AF repeater operates in the physical layer and its basic functionality is, as the name suggests, amplifying and then forwarding the received signal, including any received noise and interference.\nL2 relays operate in the data link layer and have the ability to detect and possibly correct errors that have occurred in the physical layer. L2 relays are therefore commonly called Decode-and-Forward (DF) relays as they decode the received data prior to retransmission. DF relays will, at the expense of an increased delay, not forward noise and interference.\nL3 relays operate in the network layer and are by the Third Generation Partnership Program (3GPP) regarded as being equivalent to eNodeBs (eNBs) that are wirelessly connected to a donor cell via self-backhauling. L3 relays have the same characteristics as L2 relays in the sense that they do not forward noise and interference as they perform decoding and error correction of the received signal prior to retransmission.\nThere are several known methods of utilizing various repeaters or relays to further improve the quality of performance in telecommunication systems. Some of the most common include:\nCooperative relaying which enables multiple relays to cooperate during transmissions to users. For example, the cooperation may be used for increased diversity or multiplexing of data.\nMulti-hop relaying which enables signals to be conveyed from a source to a destination over two or more wireless hops. The multiple hops are achieved by relaying the signal via one or more relay(s)/repeater(s). It may be used to reduce the end-to-end path loss and thus extending the coverage.\nOn frequency relays (or repeaters) (OFR) which are relays (repeaters) that forward on the same frequency band occupied by the received signal.\nFrequency translating relays (or repeaters) (FTR) which are relays (repeaters) that translate the retransmitted signal to another frequency band that is different from that occupied by the received signal.\nFrequency selective relays (or repeaters) (FSR) which are relays (repeaters) that may dynamically retransmit coordinated parts of the received signal bandwidth.\nThe increased path gain that comes from splitting of the signal path in two hops by repeating or relaying in an intermediate node brings several benefits: Data rates can be significantly increased; transmit power can be reduced and inter-cell interference falls rapidly. A multi-hop solution based on AF repeaters has some interesting characteristics compared to other DF relaying solutions. Since a repeater can receive and transmit on the same radio resource, which is not possible for DF types of relays, it is possible to operate without any duplex coordination loss between the two hops. A decode-and-forward (DF) relay can forward the data on the same frequency resource. However, since the decoding operation will result in an unavoidable delay the forwarding must take place at a later time instance, i.e. on another radio resource. In contrast, an amplify-and-forward (AF) repeater has a delay that typically is negligible compared to the transmission time interval, hence it can forward on the same radio resource. Repeaters also introduce less delay than DF relays which is beneficial for the performance of higher layer protocols such as TCP. Furthermore, a repeater is a simple device that typically is fairly cost efficient.\nIn particular, the use of OFR in OFDM based systems is interesting if the delay of the repeater is limited to the length of cyclic prefix of the OFDM modulation. In the air, the repeated signal path and the direct signal path add in the same way as normal multi-path does. In case of LTE, the additional time dispersion induced by the repeater does not result in any additional receiver complexity and/or reduced performance due to increased self interference as long as the total time dispersion is limited to the length of the cyclic prefix. Note that is not the case for single-carrier systems without cyclic prefix e.g. HSPA, where additional time dispersion typically increases the receiver complexity (i.e. more rake fingers are required) as well as the self-interference (i.e. signal components with a relative delay difference are non-orthogonal).\nDespite the benefits of utilizing AF repeaters, there are a few disadvantages that prevent the use from providing the full benefit of them.\nOne drawback with repeaters compared to DF relays is that they forward not only signals but also noise and interference. Furthermore, a major challenge for on-frequency repeaters is to sufficiently suppress the self-interference they induce.\nRepeaters (and relays) are efficient for both providing coverage in areas without coverage (see upper part of FIG. 1) and also to provide increased data rates to areas with weak signal strength (lower part of FIG. 1). This distinction is important since in the data-rate extension case the users will receive both a direct signal part as well as a repeated signal path, witch in the case of DF relaying will interfere with each other and in case of AF repetition and OFDM will add like multi-path. In addition, since AF repeaters amplify noise and interference they are only beneficial in case it is possible to replace one weak radio link with two significantly better radio links. This is more likely to be possible when the original radio link is weak due to some obstacle (e.g. a wall) that hinders the radio waves rather than pure propagation distance.\nIn addition, in the data rate extension scenario it is possible to dynamically turn the repeater on and off without losing coverage. It is also possible to do frequency selective repetition in the data-rate extension case without destroying the communication on the uplink and downlink control channels. Furthermore, the data-rate extension case is also particularly relevant for LTE and LTE-Advanced since, in order to compete with HSPA, their main business case is to provide high data-rates, which are only achievable in case the signal strength also is high.\nIn the coverage extension scenario, the options when it comes to advanced repeater behavior are more limited. It is not possible to e.g. turn the repeater off even when the repeater does not serve any UEs since that would leave the area with no coverage. In that case a UE wanting to perform an initial access would not be able to read the broadcast channel (BCH) and the system information blocks (SIBs) required for random access. Furthermore, an idle UE would not detect any paging messages sent from the network. In the coverage area extension case it is also not possible to perform frequency selective repetition on the downlink since that would hinder the UE from receiving the physical downlink control channel (PDCCH) that covers the whole downlink bandwidth. Also any frequency selective operation by the repeater on the uplink band must assure that the resources used for the physical uplink control channel (PUCCH) as well as the physical random access channel (PRACH) are always repeated.\nConsequently, there is a need for a more efficient use of repeaters."}
{"text": "1. Field of the Invention\nThe present invention relates to a fixing apparatus provided in an image forming apparatus, such as a copying machine, a printer, and a facsimile, and a roller used in the fixing apparatus, and in particular, a fixing apparatus of an image forming apparatus and a roller using an induction heating system.\n2. Description of the Background\nIn a fixing apparatus that adopts an induction heating system and is used in an image forming apparatus, such as an electrophotographic copying machine or printer, there is a device that increases the fixing speed by speeding up a warming up time of the fixing apparatus. For example, JP-A-2002-295452 discloses a heating device that heats a metal sleeve, which is positioned on the outer periphery of an elastic layer of a heat roller and has a small heat capacity, using an induction coil to thereby shorten a warming-up time.\nHowever, in the known device, handling of expansion or contraction that occurs to the elastic layer when warming up or cooling the elastic layer is not mentioned.\nTherefore, in a case when a metal belt having a metal layer on the outer periphery of an elastic layer is provided, it is preferable to develop a fixing apparatus of an image forming apparatus capable of increasing the life of the metal belt and an elastic roller by preventing the metal belt or the elastic layer from being broken at an early stage regardless of expansion or contraction occurring to the elastic layer."}
{"text": "This invention relates to internal combustion engines, and more particularly to plastic housings for such engines.\nSmall internal combustion engines--such as those used with lawnmowers, snowblowers, generators and the like--typically have a number of distinct components or housings for key parts of the engine. For example, the fuel tank is typically a separate component from the air filter housing. Similarly, the blower housing--which includes the air passageway through which forced air passes to cool the engine--is typically a separate component.\nEach of these distinct components must be separately manufactured and assembled, greatly increasing the cost of the completed engine. Moreover, additional parts are typically required when these housings are distinct units. For example, a fuel line is needed between the fuel tank and the carburetor. Also, a flexible hose is often used between the air filter outlet and the air intake of the carburetor or other air/fuel mixing device. The use of these additional components also increases the cost in manufacturing and assembling the completed engine.\nSince internal combustion engines by their very nature generate a great deal of heat, it has often been necessary to manufacture many of the engine components, such as the fuel tank, from metal materials due to the metal's ability to withstand the high temperatures generated by the engine. Since metal is expensive and requires a number of separate forming, punching and machining steps, the use of metal materials also increases the cost in manufacturing and assembling the component parts."}
{"text": "The present invention is directed to a conveyor system and more particularly to a roller conveyor system and various sub-assemblies thereof.\nConveyor systems are widely used within industry to transport raw materials, components and/or finished products along an assembly line or otherwise within or between manufacturing facilities. One common conveyor system is a belt driven roller conveyor system. In belt-drive roller conveyor systems, a moving belt is raised or lowered beneath a set of elongated cylinders (i.e., rollers) to make or remove contact between the moving belt and the rollers. When the moving belt contacts the rollers, the rollers rotate in the opposite direction as the belt. As a result, a bale of goods or other article situated on top of the rollers is conveyed along the conveyor path as the rollers rotate in place. Other roller conveyor systems are known, including gravity and chain driven systems.\nUnfortunately, numerous drawbacks are associated with conventional roller conveyor systems. Among those disadvantages include that the systems are often labor-intensive to install and maintain. Because they are often used continuously in a manufacturing or warehouse environment, roller conveyor systems can be subjected to long and rigorous operating conditions, resulting in wear and tear of components that require frequent maintenance. Maintenance of roller conveyor systems is often expensive, due in large part to the procurement and installation of spare parts, many of which are heavy and cumbersome.\nStill other drawbacks not adequately addressed with known roller conveyor systems relate to safety. For example, individuals cannot cross a roller conveyor system without risk of falling onto the rollers, which can spin freely underfoot as one tries to walk across them. However, the alternative, walking entirely around the roller conveyor system, is not efficient, particularly since roller conveyor systems can extend for tens of yards or more.\nIt would be desirable in the art for a roller conveyor system without the above-mentioned drawbacks."}
{"text": "The present invention generally relates to full color liquid crystal displays generated by the use of birefringent filters. More particularly, the present invention relates to full color liquid crystal displays using polarization interference filters in a passive matrix addressing scheme.\nTraditionally color generation in liquid crystal displays (LCD) is accomplished by the patterning of color filters onto individual pixels. However, the use of color filters results in both a loss of resolution, due to subpixeling, and lower light transmission, due to the absorptive nature of the filters. Other methods of generating color include field sequential color in which the actual color is generated by RGB light emitting diodes (LED) and the liquid crystal (LC) pixel is driven so as to either block or transmit the colored light. This typically requires the LC to be driven at a very high frame rate which in turn calls for a very fast switching LC material and mode. Such materials and modes are difficult to find.\nAnother method of color generation employs polarization interference filters (PIF). Such filters work by introducing a phase shift between two orthogonally polarized field components by either a static or dynamic optical element such as a uniaxial retarder or electrically driven LC cell. Color is generated by the interference of these two components with an analyzer, and color switching is accomplished by changing the phase shift between the two components by using dynamic optical elements. Since there is no more absorption than the losses associated with polarizer absorption, the PIF can generate color with a higher luminance than absorptive color filters.\nThe use of retarders between polarizers to function as filters has been studied for several years, starting with the original design for monochromatic imaging by French astronomer Lyot. See I. J. Hodgkinson, Q. H. Wu: Birefringent Thin Films and Polarizing Elements, (World Scientific Press, Singapore 1997) 1st ed., Chap. 4, p. 52, incorporated herein by reference. The Lyot filter design was improved upon by the Solc design, which offered similar filtering characteristics but higher overall transmission. This design is generally disclosed in Ivan Solc, Czechoslovak Journal of Physics, Vol. 10(1), 16-34 (1960), also incorporated herein by reference. The Solc design was then generalized by Harris et al., A Optical Network Synthesis Using Birefringent Crystals, 1. Synthesis of lossless Networks of Equal-Length Crystals, J. Opt Soc. America (1964) 54(10):1267-1279, incorporated herein, describing a synthesis procedure for obtaining various filter designs.\nThe two original Solc designs are the “fan” and “folded” type. In the fan design, the wavelengths at which the elements have an even number of half waves of retardation pass through an output polarizer oriented parallel to the input polarizer. In the folded design, the wavelengths at which the elements have an odd number of half waves of retardation pass through the output polarizer, which is crossed with respect to the input polarizer. The passband characteristics of the Solc filters are directly related to the number of elements used to create the filter, i.e., the larger the number of elements, the sharper the passband. Another crucial element in the transmission function is the size and frequency of the sidelobes. The Solc filters, as it turns out, are not optimal in design and give passband sidelobes that can be avoided using various synthesis techniques.\nThe use of PIFs for the generation of color in a display device was shown by Sharp et al., U.S. Pat. No. 5,929,946. Starting from a Solc filter design, the design of a device capable of switching between a single desired color and black or white, depending upon polarizer orientations, was shown. Given that a folded Solc filter will filter the desired wavelengths by rotating them by pi/2, a second active optical element was added to the filter in order to be able to switch between a filtering and non-filtering action. Particularly, a single, switchable half wave plate was placed between the filter and the second half of the filter was repositioned. Proper repositioning required inverting the second half of the filter and rotating it by pi/2, as shown in FIG. 1, and discussed below.\nWith reference to FIG. 1, a single color generation stage of a display device according to Sharp is designated generally by the numeral 1. Color generation stage 1 includes a first retarder stack 2 (pre filter), a second retarder stack 3 (post filter), which is inverted and rotated by 90° with respect to first retarder stack 2, and an active matrix addressed optical element 4 (liquid crystal cell), particularly, a half wave plate, placed between first and second retarder stacks 2,3. First retarder stack 2 is oriented to rotate a desired wave length of light to an angle of pi/4, with respect to the optic axis of the active matrix addressed optical element 4 (half wave plate). As is generally known, crossed or parallel polarizers (not shown), particularly an input polarizer and an output analyzer, are employed to introduce polarized light to the color generation stage 1 and to analyze the light transmitted through color generation stage 1.\nIf LC cell 4 is switched so that its net retardation is zero, then the first part of the filter, first retarder stack 2, is crossed with respect to the second part of the filter, second retarder stack 3, as shown in FIG. 1. In this case, the net retardation between the two polarizers is zero and so the display appears black for crossed polarizers and white for parallel polarizers. If LC cell 4 is switched so that it is a half wave plate, then its effect on the filter depends on its orientation in the stack. Given that a folded Solc design rotates the desired wavelength by pi/2, then half the filter will rotate it by pi/4. As shown in FIG. 1, light 5 entering the filter from the left will have the desired wavelength rotated by pi/4 while light entering from the right will be rotated by −pi/4 to 3pi/4. Therefore, for light 5 to pass through the complete filter, LC cell 4 must take light polarized at pi/4 and rotate it to 3pi/4 degrees. A pi/2 rotation is easily accomplished by placing LC cell 4 at pi/4 to the incoming polarization. In this position the LC cell allows the filter to be switched on or off depending on its own state. For intermediate voltages, LC cell 4 acts as a retarder between 0 and pi/2 and so rotates the incoming linearly polarized light to some elliptical state. As a result, the desired wavelength is not completely rotated by the second half of the filter and so suffers absorptive losses at the analyzer. In this manner, intermediate grey scale colors can be created.\nThe effect of the filter on the rest of the spectra is as follows. In the case of a pure folded Solc filter, the undesired wavelengths entering the filter linearly polarized at 0° are all left unrotated (more so for a multiple plate retarder than a two plate) and are absorbed at the analyzer, which is crossed with respect to the polarizer. For the above design, the placement of the LC cell ensures that the unwanted spectra remains unaffected. While the desired wavelength is rotated, by the first part of the filter, to pi/4, the rest of the spectra remains at 0°. Thus, the undesired wavelengths are at pi/2 to the LC cell optics axis and the LC cell has no effect on them. Passing through the second half of the filter, they remain linearly polarized at 0° and are absorbed by the analyzer.\nTherefore, in summary, Sharp design requires the following:                (1) a polarizer providing linearly polarized white light at a known orientation;        (2) a prefilter to rotate desired wavelengths to an angle of pi/4 with respect to the optic axis of the LC cell (or other electro- or magneto-optic modulator) while leaving undesired wavelengths untouched;        (3) an active matrix addressed optical element such as a LC cell (or other electro- or magneto-optic modulator) that rotates linearly polarized light by pi/2, i.e., a half wave plate, and has no effect on light polarized along or orthogonal to its optic axis; and        (4) a postfilter that is a copy of the prefilter except inverted and rotated by pi/2.        \nSuch a filter design is capable of producing a relatively narrow passband transmission function (depending on the number of retarders making up the pre and post filter) and so can produce a single color or its grey scales from input white light. Such a filter is discussed in S. Saeed, P. J. Bos, Z. Li, SID00 Digest, 830 (2000), incorporated herein by reference. In order to produce more than one color, i.e., red, green, and blue, a stacked approach is taken to provide a full color liquid crystal display, as shown in FIGS. 2 and 3. In the stacked approach, a color generation stage according to FIG. 1 is tuned for each of the three primary colors (designated as 1B, 1G, and 1R) and then stacked one on top of the other. In FIG. 2, like parts as compared to FIG. 1 have received like numerals, with distinctions made between each stage and its elements by employing the letters B, G, and R to indicate blue, green, and red. Only one polarizer 6 and analyzer 7 is used at the two ends, and, in FIGS. 2 and 3, they have transmission axes that are orthogonal to one another, although parallel configuration may alternately be practiced. The stacked design works since each stack only rotates the wavelengths associated with its color, eg., the blue stack is designed to rotate wavelengths 430 nm to 490 nm while the rest of the spectra is left linearly polarized at 0°. Since the analyzer is placed at the end of the stack, this part of the spectra is able to move through the other two stacks, as shown in FIG. 3(B).\nIn FIG. 3, white polarized light is shown passing through three different stacks of filters, each stack containing three filters individually tuned for red (1R), green (1G), and blue (1B) wavelengths. At (A), the filters are all in an optically inactive mode, and the entire spectrum is blocked by the crossed polarizers. At (B), the blue filter is in an optically active mode, such that the blue spectrum is rotated and the stack ultimately outputs blue light. At (C), all three of the filters, blue (1B), green (1G), and red (1R), are in an optically active mode, and, thus, all three parts of the spectrum, blue, green, and red, are rotated, and the stack ultimately outputs white light.\nSince the entire white light spectra is able to move through the three layer stack with each layer only affecting selected wavelengths, the stack can produce red, green, blue and white light or any combination thereof. This however will only be true if the red, green, and blue spectra do not overlap. As mentioned before, the bandwidth of each filter depends on the number of elements making up the filter. If only four retarders are used, for example, then the filter's bandwidth will be quite broad. If this filter was to be placed in the stacked design, it would rotate parts of the undesired spectra along with the desired spectra. For example, if the red filter is broadband, then it will rotate parts of the green spectra also. As a result, if the stack is setup to be in the white state, i.e., all three layers are rotating their selected wavelengths, then the green would be over rotated, and the resulting white output would lack some green. Therefore, in a stack design, it is important to have each filter's bandpass be as narrow as possible so as to prevent the overlap that reduces efficiency.\nThere are a number of orientations that the retarders making up the filter can have and still provide similar if not identical transfer functions. Synthesis techniques described by Harris et al., cited above, and Buhrer, Appl. Opt. 33, 2249-2254 (1994), are useful for obtaining the best orientations of the filter elements for a required transfer function, and each reference cited is incorporated herein. These synthesis techniques rely on the fact that each element in the filter provides one cosine Fourier component of the desired transmission function. Such synthesis techniques are useful in situations where the active element in the filter operates in the 0 to pi/2 range.\nA birefringent filter based LCD employing electrically controlled birefringence (ECB) type cells, which fall into the category of optically simple active elements, was disclosed in the above-referenced Saeed et al. publication. The ECB type cells, however, require a rather expensive drive method called active matrix addressing, which involves the use of a transistor at each pixel edge to provide a switching operation.\nA less expensive drive method is passive matrix addressing. A passive matrix addressing design could make use of the electro-optical characteristics of the LC mode and could drive displays without the use of transistors, and, as a result, the display would be cheaper and easier to manufacture. However, the requirement of a sharp electro-optical curve means that LC modes such as the ECB cannot be used in conjunction with passive matrix addressing. Thus, there exists a need in the art for a birefringent filter-based color generation scheme for a passive matrix display device.\nIn the design of a single stage of a color filter, the incident light can be considered to have two spectral components: a controlled spectral component and an un-controlled spectral component.\nOne of the required features of the electro-optical element in the color filter of Sharp is that the polarization state of the controlled spectral component of the light is effected by the state of the electro-optical element, but the polarization state of the uncontrolled component is not.\nSharp used an untwisted liquid crystal device for which it is well known that if light is linearly polarized with a particular state (at 45° to the projection of the optic axis on the plane of the device), that the polarization state of the light will be affected by a voltage applied to the device; but if light is polarized parallel to the projection of the optic axis, then the polarization state will not be affected by a voltage applied to the device.\nHowever, untwisted nematic devices typically do not show that appropriate voltage versus polarization state change that is required for a multiplexed passive matrix display, and it is not clear, from Sharp, that it could be possible to use a nematic material of twisted structure. In fact, it is well known that twisted devices do not have the characteristic that a state of linearly polarized light with a wide wavelength spectral band can be unaffected by the state of the device. Therefore, the problem of designing a filter of the type discussed by Sharp, but having a twisted, multiplexable, liquid crystal device, is the focus of this invention."}
{"text": "1. Field of the Invention\nThis invention relates to computer display systems and, more particularly, to methods and apparatus for providing a graphics accelerator capable of selectively providing during any clock period color values for at least two pixels blended with a single texture or color values for a single pixel blended with a plurality of textures.\n2. History of the Prior Art\nIn three dimensional graphics, surfaces are typically rendered by assembling a plurality of polygons in a desired shape. The polygons are conventionally triangles having vertices which are defined by three dimensional coordinates in world space, by color values, and by texture coordinates.\nTo display a surface on a computer monitor, the three dimensional world space coordinates are transformed into screen coordinates in which horizontal and vertical values (x, y) define screen position and a depth value (z) determines how near a vertex is to the screen and thus whether that vertex is viewed with respect to other points at the same screen coordinates. The color values (r, g, b) define the brightness of each of red/green/blue colors at each vertex and thus the color (often called diffuse color) at each vertex. Texture coordinates (u, v) define texture map coordinates for each vertex on a particular texture map defined by values stored in memory.\nA texture map typically describes a pattern to be applied to the surface of the triangle to vary the color in accordance with the pattern. The texture coordinates of the vertices of a triangular surface area fix the position of the vertices of the triangle on the texture map and thereby determine the texture detail applied to each portion of the surface within the triangle in accordance with the particular texture. In turn, the three dimensional coordinates of the vertices of a triangle positioned on the texture map define the plane in which the texture map and the surface lie with respect to the screen surface.\nA texture which is applied to a surface in space may have a wide variety of characteristics. A texture may define a pattern such as a stone wall. It may define light reflected from positions on the surface. It may describe the degree of transparency of a surface and thus how other objects are seen through the surface. A texture may provide characteristics such a dirt and scratches which make a surface appear more realistic. A number of other variations may be provided which fall within the general description of a texture. In theory, a number of different textures may be applied to any triangular surface.\nIn order to apply more than one texture to a surface, prior art graphics accelerators initially were designed to progress through a series of steps [is] in which pixel coordinates and color values describing each triangle are first generated one pixel at a time in sequence, a first texture is mapped to the triangle using the texture coordinates of the vertices and texture coordinates are generated for each pixel as the pixel coordinates are generated, texture values describing the variation of each pixel in the triangle for the first texture are generated using the texture coordinates for each pixel, the texture value describing the first texture for one pixel and the diffuse color values describing that pixel of the triangle are blended to produce a textured color value for the pixel, and the resulting triangle generated from all of the textured color values is blended with any image residing in a frame buffer from which the image may be presented on an output display. Then, texture values for a second texture mapped to the same triangle are generated and blended with the same sequence of pixel color values in the same manner, and the triangle blended with the second texture is blended with the image residing in the frame buffer.\nThe need to transit the graphics pipeline to blend each texture to the surface of each triangle defining an output image slows the process drastically. In many cases involving rapidly changing images, it has limited significantly the number of textures which can be applied. For this reason, a more recent development provides a pair of texture stages and a pair of texture blend stages in the pipeline. The first texture stage generates texture values describing a first texture from texture coordinates provided as each pixel is generated. These first texture values are blended with the pixel color values at the first texture blend stage as each set of pixel color values is generated. The textured color value output of the first texture blend stage is then furnished to the second texture blend stage. The textured color value output of the first texture blend stage is blended with texture values generated by the second texture stage using texture coordinates provided as each pixel is generated. The output of the second texture blend stage is ultimately transferred to the frame buffer blend stage to be blended with the image data already in the frame buffer. This more advanced pipeline allows two textures to be blended with a surface in a single pass through the graphics pipeline.\nAlthough this most recent development is useful in accelerating texture blending in a graphics pipeline, it is limited to producing a single pixel having at most two textures during any clock of the graphics pipeline and cannot be utilized for any other purposes. More complicated functions require the use of the host processor and the frame buffer blending stage and drastically slow the rendering of surfaces by the graphics accelerator.\nThere are no prior art systems which have been capable of providing two or more textured pixels during each clock period.\nIt is desirable to provide a new computer graphics pipeline capable of rapidly selectively providing a plurality of textured pixels or a single pixel blended with a plurality of textures during any clock period.\nThe present invention is realized by a graphics accelerator pipeline including a rasterizer stage, a texture stage, and a combiner stage capable of producing output values during each clock interval of the pipeline which map an individual texture to a plurality of pixels."}
{"text": "The present invention relates to systems and methods for estimating switching supply load current.\nThe capability to dynamically measure the power consumption of an electronic system is highly desirable. Some of the benefits include: system fault detection (should power exceed a normal range), engineering power savings (during system prototype development), ability to provide accurate battery life estimates (in a battery operated system). Today's modern electronic systems often utilize switching regulators in order to improve the system power efficiency and reduce heat. In applications where large amounts of current are required, external power field effect transistors (FETs) (as opposed to power FETs integrated with the switching regulator) are most economical. One aspect of power estimation requires measuring the output current supplied by these external FETs.\nPrevious solutions to measuring the output current flowing in external FETs, utilize a sense resistor in series with the output, wherein the voltage drop across the resistor was proportional to the output current. In this circumstance, the output current is measured by sensing this resistor voltage drop. This method works at the expense of lower efficiency, as any voltage across the sense resistor results in a power loss and resulting degradation in the regulator efficiency. Consequently, in order to minimize the power loss, the sense resistor value is kept small which results in poor resolution of low current measurements (due to the low resulting voltage drop across the sense resistor). Also, low value resistors are expensive.\nAn alternate method for measuring the output current requires the use of specialized output inductors. Most switching regulators utilize external inductors as part of the voltage transformation/regulation loop. It is possible to add a separate set of “turns” around the output inductor which sense the magnitude of the magnetic flux in the inductor. The magnitude of the magnetic flux is proportional to the current in the inductor. Thus, the extra turns provide a means to sense the current flow in the output inductor. While this method does not suffer the efficiency loss of the sense resistor method, it does require the use of a more expensive and non-standard output inductor.\nFIG. 1 shows a typical BUCK regulator output power stage. In this regulator design, supply voltage 6 provides power to output power FET Q2 16, which is controlled by “Break Before Make” (BBM) unit 12. Output FET Q2 16 is also connected to output FET Q1 20. Body diodes D1 22 and D2 18 are built into the output power FETs Q1 20 and Q2 16. “FET Control” circuit 10 (shown as a voltage source) generates a switching signal to command the power FETs to turn on and off. BBM units 12-14 insure that under no circumstances are both power FETs 16 and 20 on at the same time, which would cause a large “shoot though” current to flow in the FETs 16 and 20. Thus, for example, if Q2 16 is on and Q1 20 are off then the BBMs 12 and 14 insure that Q2 16 is shut off before Q1 20 is turned on. Because of the BBMs 12 and 14, there is a short period of time, perhaps 20 ns or more, (called the “dead time”) when neither Q1 20 nor Q2 16 is turned on.\nThe physics of inductor 26 do not allow for instantaneous changes in the current flowing through the inductor. Thus if the power output stage is sourcing current to the load Rload 30 and capacitor C1 28, then even during the short dead time, when both Q1 20 and Q2 16 are off, the load current will continue to flow through inductor 26. During this period of time, the body diode D1 22 will source the current from ground through inductor 26 to the load. When body diode D1 22 is conducting current, the voltage at the PowerSwitch_out 24 will go below ground."}
{"text": "Technical Field\nThis disclosure relates generally to transaction processing using highly-scalable, decentralized peer-to-peer (P2P) networks, known as blockchains.\nBackground of the Related Art\nBlockchain is a type of business transaction ledger. A blockchain network is a decentralized system for the exchange of assets. It uses a shared ledger (that is either public or private) for recording the history of electronic business transactions that take place in a peer-to-peer (P2P) business network. A blockchain network uses a decentralized consensus mechanism or protocol. In particular, all validating nodes in the network run the same (agreed-upon) consensus algorithm against the same transactions, and thus validate (or invalidate) each transaction. Valid transactions update the ledger. A blockchain peer-to-peer network is resilient given its decentralized topology. As member nodes join or leave the network dynamically, messages are exchanged between the network participants on a best-effort broadcast basis.\nA blockchain is a permanent digitized chain of transactions, grouped into blocks, that ensures that participants cannot tamper with or deny past transactions. A permissioned blockchain is one in which the participants who invoke business transactions, as well as those who control and manage copies of the blockchain, are known. To this end, an identity service in a permissioned blockchain typically links information about business entities to the cryptographic keys used to digitally sign transactions. Typically, a ledger comprises a blockchain, and an associated world state. The world state typically is a view, but can be a superset, of the blockchain that comprises a set of referenceable variables that, while necessarily recorded in the blockchain, can be utilized by programs running on the network. In particular, to update the ledger with a proposed transaction, each node must process the transaction against the logic of one or more programs running directly on the replicated network. These programs, called chaincode, are replicated across each node in the network to guarantee uptime and functionality. Chaincode is stored on the blockchain and triggered by transaction requests.\nHyperledger is a cross-industry collaborative effort to support permissioned blockchain-based distributed ledgers. It is focused on ledgers designed to support global business transactions, including major technological, financial, and supply chain companies, with the goal of improving many aspects of performance and reliability. The project aims to bring together a number of independent efforts to develop open protocols and standards by providing a modular framework that supports different components for different uses. A Hyperledger is a peer-to-peer distributed network that uses cryptographic security, decentralized consensus, and a shared public ledger with properly-controlled and permissioned visibility.\nTo bring new validating peers into a blockchain network, such as a permissioned blockchain like Hyperledger, it is necessary to transmit the blockchain state from one or more of the network peers to the new peer. Long-lived blockchains will contain an extremely large number of blocks and world state. It is both inefficient, and impractical, to transmit tens or hundreds of gigabytes, if not terabytes, to the new peers in the network simply to enable these new peer nodes to participate in the network. To support the efficient and rapid addition of new validating (and non-validating) peer nodes in a blockchain network, there needs to be an efficient mechanism for identifying the current state of the blockchain so it can be efficiently communicated to the new peer nodes, and enable them to start participating in the network.\nIn addition, many types of businesses, including finance, have auditing requirements that are built on the premise that the audit will be conducted over an interval of time—from a start date to an end date. It is assumed that the state of the ledger (e.g., a financial ledger) is consistent and acceptable prior to the starting date and that the transactions from that date until the ending date are to be audited. Earlier data may be available for reference purposes, but it may not be required. As such, it be desirable to provide a way to record of the state of the system at the starting point of the audit interval and perhaps at the end of the interval as well.\nMany kinds of blockchains applications are currently under discussion and their design are intended to be very long-lived (e.g., even up to 50 years or more). For long-lived blockchain ledgers, typically there are requirements for pruning the ledger for practical or regulatory reasons. In the absence of any application-specific audit checkpoint logic, to perform an audit using the current Hyperledger design, the auditor will need to start with the ledger's genesis (first) block and run all (relevant) transactions forward through time to reach the start of the audit period, and then continue running all of the relevant transactions for the current audit period. If the ledger contains substantial history (5 years, 10 years, or longer, with millions or billions of transactions), this process of processing the transactions to reach the start of the audit period can require a substantial amount of computational resources. Thus, it would be highly desirable to find a way to forestall this lengthy and resource-intensive process.\nIn financial and other business scenarios, old data prior to specific points in time are no longer relevant and can be discarded. Some businesses have record retention policies. Data older than a particular date is to be discarded. For privacy and legal reasons, old data is to be discarded. For Internet of Things (IoT)-based applications, many desired use cases have limited storage capacity, and thus pruning is essential. In such scenarios, bringing a new permissioned blockchain validating peer online requires transfer of the ledger (blockchain, world state) to the new peer, and efficiently enabling this operation may be critical."}
{"text": "Beverage preparation machines have been known for a number of years. For example, U.S. Pat. No. 5,943,472 discloses a water circulation system between a water reservoir and a hot water or vapour distribution chamber of an espresso machine. The circulation system includes a valve, metallic heating tube and pump that are connected together and to the reservoir.\nA beverage preparation machine typically includes a housing containing a beverage processing module and a water tank in fluid communication with the beverage processing module. Examples of such beverage preparation machines are disclosed in EP 1 208 782, EP 1 267 687, EP 1 686 879, EP 1 731 065, EP 1 829 469, EP 1 864 598, EP 1 865 815, EP 1 867 260, EP 1 878 368, EP 2 222 210, EP 2 222 211, EP 2 222 212, EP 2 227 121, EP 2 227 122, US 2008/0006159, U.S. Pat. No. 7,165,488, WO 2007/111884, WO 2009/074553, WO 2010/015427 and WO 2012/055767.\nUsually the water tank is removable to be refilled by a user when empty. Some systems include a continuous water supply by connecting the beverage machine directly to the city water distribution network, as for instance disclosed in CN201076369. This publication discloses a system that has a refillable water tank as well as an inlet for direct connection with the tap. Hence, beverage preparation can be carried out either by using water directly from the tap or from the water tank."}
{"text": "Wireless telecommunication devices are evolving to contain increased functionality and complexity. This increased functionality often brings together functions that have been traditionally been provided by different devices such as cell phones and personal digital assistants (PDA). The combination of these functions typically require increased processor capability as well as increased power requirements. The requirement for additional processing capability to add functionality and minimize latency is especially important when the information must be processed in real time, as for example in cell phones. Having more processing capability and in turn higher power consumption, is especially problematic in wireless communication systems where it is inconvenient to connect to power sources.\nWireless communication systems generally must contain their own source of power, which often is in the form of a battery. Users typically need the ability to operate such systems for longer periods of time without the need to recharge or swap batteries or even connect to line power. However, such longer operating times normally require an increase in battery size, which leads to undesirable effects such as heavier batteries, increased expense, and environmental concerns regarding disposal of used batteries.\nTo meet the needs of increased processing power within wireless communication devices, additional processors requiring more memory and power were added to devices. A general purpose processor handles most system tasks and a modem computing subsystem handles tasks related to handling mobile station requirements. Mobile station modem binary software images (i.e., contents of memory) are programmed at the time of manufacture into read-only memory (ROM) as a single contiguous binary image. The modem computing subsystem directly executes the memory image from ROM, which results in slower execution than images executed from memories with faster access times, such as random access memory (RAM). At system boot time (e.g., when the wireless device is powered-up), read-write and zero-initialized data are copied to RAM prior to execution of code by the modem computing subsystem. No part of the over-the-air standards as implemented in the software binary image can be executed prior to completion of system boot and initiation of the operating system. All system memory must be completely powered-up prior to mobile station modem operation of the over-the-air standard. This approach results in wasting significant amounts of power because the modem computer subsystem had to be powered up even when not in use.\nTherefore, it can be appreciated that there is a significant need for a system and method to minimize power consumption in a wireless communication system while increasing the speed and functionality of the device for the user. The present invention provides this and other advantages that will be apparent from the following detailed description and accompanying figures."}
{"text": "1. Field of the Invention\nThe invention relates generally to ultra-wideband communications, and more particularly to systems and methods for communication using ultra-wideband technology.\n2. Background\nWireless communication systems are proliferating at the Wide Area Network (WAN), Local Area Network (LAN), and Personal Area Network (PAN) levels. These wireless communication systems use a variety of techniques to allow simultaneous access to multiple users. The most common of these techniques are Frequency Division Multiple Access (FDMA), which assigns specific frequencies to each user, Time Division Multiple Access (TDMA), which assigns particular time slots to each user, and Code Division Multiple Access (CDMA), which assigns specific codes to each user. But these wireless communication systems and various modulation techniques are afflicted by a host of problems that limit the capacity and the quality of service provided to the users. The following paragraphs briefly describe a few of these problems for the purpose of illustration.\nOne problem that can exist in a wireless communication system is multipath interference. Multipath interference, or multipath, occurs because some of the energy in a transmitted wireless signal bounces off of obstacles, such as buildings or mountains, as it travels from source to destination. The obstacles in effect create reflections of the transmitted signal and the more obstacles there are, the more reflections they generate. The reflections then travel along their own transmission paths to the destination (or receiver). The reflections will contain the same information as the original signal; however, because of the differing transmission path lengths, the reflected signals will be out of phase with the original signal. As a result, they will often combine destructively with the original signal in the receiver. This is referred to as fading. To combat fading, current systems typically try to estimate the multipath effects and then compensate for them in the receiver using an equalizer. In practice, however, it is very difficult to achieve effective multipath compensation.\nA second problem that can affect the operation of wireless communication systems is interference from adjacent communication cells within the system. In FDMA/TDMA systems, this type of interference is prevent through a frequency reuse plan. Under a frequency reuse plan, available communication frequencies are allocated to communication cells within the communication system such that the same frequency will not be used in adjacent cells. Essentially, the available frequencies are split into groups. The number of groups is termed the reuse factor. Then the communication cells are grouped into clusters, each cluster containing the same number of cells as there are frequency groups. Each frequency group is then assigned to a cell in each cluster. Thus, if a frequency reuse factor of 7 is used, for example, then a particular communication frequency will be used only once in every seven communication cells. Thus, in any group of seven communication cells, each cell can only use 1/7th of the available frequencies, i.e., each cell is only able to use 1/7th of the available bandwidth.\nIn a CDMA communication system, each cell uses the same wideband communication channel. In order to avoid interference with adjacent cells, each communication cell uses a particular set of spread spectrum codes to differentiate communications within the cell from those originating outside of the cell. Thus, CDMA systems preserve the bandwidth in the sense that they avoid reuse planning. But as will be discussed, there are other issues that limit the bandwidth in CDMA systems as well.\nThus, in overcoming interference, system bandwidth is often sacrificed. Bandwidth is becoming a very valuable commodity as wireless communication systems continue to expand by adding more and more users. Therefore, trading off bandwidth for system performance is a costly, albeit necessary, proposition that is inherent in all wireless communication systems.\nThe foregoing are just two examples of the types of problems that can affect conventional wireless communication systems. The examples also illustrate that there are many aspects of wireless communication system performance that can be improved through systems and methods that, for example, reduce interference, increase bandwidth, or both.\nNot only are conventional wireless communication systems effected by problems, such as those described in the preceding paragraphs, but also different types of systems are effected in different ways and to different degrees. Wireless communication systems can be split into three types: 1) line-of-sight systems, which can include point-to-point or point-to-multipoint systems; 2) indoor non-line of sight systems; and 3) outdoor systems such as wireless WANs. Line-of-sight systems are least affected by the problems described above, while indoor systems are more affected, due for example to signals bouncing off of building walls. Outdoor systems are by far the most affected of the three systems. Because these types of problems are limiting factors in the design of wireless transmitters and receivers, such designs must be tailored to the specific types of system in which it will operate. In practice, each type of system implements unique communication standards that address the issues unique to the particular type of system. Even if an indoor system used the same communication protocols and modulation techniques as an outdoor system, for example, the receiver designs would still be different because multipath and other problems are unique to a given type of system and must be addressed with unique solutions. This would not necessarily be the case if cost efficient and effective methodologies can be developed to combat such problems as described above that build in programmability so that a device can be reconfigured for different types of systems and still maintain superior performance."}
{"text": "Aptamers, are oligonucleotides, which can be synthetic or natural, that bind to a particular target molecule, such as a protein or metabolite. Typically, the binding is through interactions other than classic Watson-Crick base pairing.\nAptamers represent a promising class of therapeutic agents currently in pre-clinical and clinical development. Like biologics, e.g., peptides or monoclonal antibodies, aptamers are capable of binding specifically to molecular targets and, through binding, inhibiting target function. A typical aptamer is 10-15 kDa in size (i.e., 30-45 nucleotides), binds its target with sub-nanomolar affinity, and discriminates among closely related targets (e.g., will typically not bind other proteins from the same gene family) (Griffin, et al. (1993), Gene 137(1): 25-31; Jenison, et al. (1998), Antisense Nucleic Acid Drug Dev. 8(4): 265-79; Bell, et al. (1999), In Vitro Cell. Dev. Biol. Anim. 35(9): 533-42; Watson, et al. (2000), Antisense Nucleic Acid Drug Dev. 10(2): 63-75; Daniels, et al. (2002), Anal. Biochem. 305(2): 214-26; Chen, et al. (2003), Proc. Natl. Acad. Sci. U.S.A. 100(16): 9226-31; Khati, et al. (2003), J. Virol. 77(23): 12692-8; Vaish, et al. (2003), Biochemistry 42(29): 8842-51).\nAptamers can be created by an entirely in vitro selection process (Systematic Evaluation of Ligands by Experimental Enrichment, i.e., SELEX™) from libraries of random sequence oligonucleotides as described in U.S. Pat. Nos. 5,475,096 and 5,270,163. Aptamers have been generated against numerous proteins of therapeutic interest, including growth factors, enzymes, immunoglobulins, and receptors (Ellington and Szostak (1990), Nature 346(6287): 818-22; Tuerk and Gold (1990), Science 249(4968): 505-510).\nAptamers have a number of attractive characteristics for use as therapeutics. In addition to high target affinity and specificity, aptamers have shown little or no toxicity or immunogenicity in standard assays (Wlotzka, et al. (2002), Proc. Natl. Acad. Sci. U.S.A. 99(13): 8898-902). Indeed, several therapeutic aptamers have been optimized and advanced through varying stages of pre-clinical development, including pharmacokinetic analysis, characterization of biological efficacy in cellular and animal disease models, and preliminary safety pharmacology assessment (Reyderman and Stavchansky (1998), Pharmaceutical Research 15(6): 904-10; Tucker et al., (1999), J. Chromatography B. 732: 203-212; Watson, et al. (2000), Antisense Nucleic Acid Drug Dev. 10(2): 63-75).\nIt is important that the pharmacokinetic properties for all oligonucleotide-based therapeutics, including aptamers, be tailored to match the desired pharmaceutical application. While aptamers directed against extracellular targets do not suffer from difficulties associated with intracellular delivery (as is the case with antisense and RNAi-based therapeutics), the aptamer must be distributed to target organs and tissues, and remain in the body (unmodified) for a period of time consistent with the desired dosing regimen. Early work on nucleic acid-based therapeutics has shown that, while unmodified oligonucleotides are degraded rapidly by nuclease digestion, protective modifications at the 2′-position of the sugar, and use of inverted terminal cap structures, e.g., [3′-3′dT], dramatically improve nucleic acid stability in vitro and in vivo (Green, et al. (1995), Chem. Biol. 2(10): 683-95; Jellinek, et al. (1995), Biochemistry 34(36): 11363-72; Ruckman, et al. (1998), J. Biol. Chem. 273(32): 20556-67; Uhlmann, et al. (2000), Methods Enzymol. 313: 268-84). In some SELEX selections (i.e., SELEX experiments or SELEX ions), starting pools of nucleic acids from which aptamers are selected are typically pre-stabilized by chemical modification, for example by incorporation of 2′-fluoropyrimidine (2′-F) substituted nucleotides, to enhance resistance of aptamers against nuclease attack. Aptamers incorporating 2′-O-methylpurine (2′-OMe purine) substituted nucleotides have also been developed through post-SELEX modification steps or, more recently, by enabling synthesis of 2′-OMe-containing random sequence libraries as an integral component of the SELEX process itself.\nIn addition to clearance by nucleases, oligonucleotide therapeutics are subject to elimination via renal filtration. As such, a nuclease-resistant oligonucleotide administered intravenously exhibits an in vivo half-life of <10 min, unless filtration can be blocked. This can be accomplished by either facilitating rapid distribution out of the blood stream into tissues or by increasing the apparent molecular weight of the oligonucleotide above the effective size cut-off for the glomerulus. Conjugation to a PEG polymer (“PEGylation”) can dramatically lengthen residence times of aptamers in circulation, thereby decreasing dosing frequency and enhancing effectiveness against targets. Previous work in animals has examined the plasma pharmacokinetic properties of PEG-conjugated aptamers (Reyderman and Stavchansky (1998), Pharmaceutical Research 15(6): 904-10; Watson, et al. (2000), Antisense Nucleic Acid Drug Dev. 10(2): 63-75)). Determining the extravasation of an aptamer therapeutic, including aptamer therapeutics conjugated to a modifying moiety or containing modified nucleotides and, in particular, determining the potential of aptamers or their modified forms to access diseased tissues (for example, sites of inflammation, or the interior of tumors) define the spectrum of therapeutic opportunities for aptamer intervention.\nTypically, therapeutic aptamers are administered by injection, for example, by subcutaneous injection. Accordingly, the aptamer must be dissolved in a liquid vehicle for administration. The relatively high molecular weight of aptamers, and in particular aptamers that have been derivatized, for example by PEGylation, however, often makes it difficult to obtain a pharmaceutical composition wherein the aptamer is dissolved in a pharmaceutically acceptable solvent at a sufficient concentration to provide a pharmaceutical composition that is clinically useful for administration to an animal.\nU.S. published application no. 2005/0175708 discloses a composition of matter that permits the sustained delivery of aptamers to a mammal. The aptamers are administered as microspheres that permit sustained release of the aptamers to the site of interest so that the aptamers can exert their biological activity over a prolonged period of time. The aptamers, can be anti-VEGF aptamers.\nP. Burmeister et al., (2004), Chemistry and Biology: 15, 25-33 disclose a method for generating a 2′-O-methyl aptamer (ARC245) that binds to vascular endothelial growth factor, which exhibits good stability.\nAccordingly, there is a need in the art for improved pharmaceutical compositions, wherein the therapeutic agent is an aptamer. In particular, there is a need for pharmaceutical composition wherein the aptamer can be dissolved in a pharmaceutically acceptable solvent at a sufficient concentration to provide a pharmaceutical composition that is clinically useful for administration to an animal. The present invention addresses this as well as other needs.\nCitation of any reference in this application is not to be construed as an admission that such reference is prior art to the present application."}
{"text": "This invention relates generally to control systems, and more particularly to a controller for controlling flow, speed, pressure or performance of a pumping system.\nA typical centrifugal pump of the prior art comprises an impeller, rotatably mounted in a stationary casing with the rotating impeller imparting pressure and kinetic energy to the fluid being pumped, and the stationary casing guiding the fluid to and from the impeller. In a typical centrifugal pump casing, which generally includes concentric, diffusor and, volute type centrifugal casings, the rotation of the impeller imparts kinetic energy to the fluid and causes fluid flow, in a generally circular direction about the perimeter of the impeller, through the casing surrounding the impeller. At some point in the casing, the fluid flows from the perimeter of the impeller, passes a cut-water or the like through an area of the pump generally known as the discharge inlet area and through the discharge nozzle to the pump discharge.\nThe fluid flow can be affected by the design of the impeller, the design and size of the casing, the speed at which the impeller rotates, and design and size of the pump inlet and outlet, quality and finish of the components, presence of a casing volute and the like. In order to control fluid flow, variable frequency devices have been used to adjust the motor speed of the pump so as to regulate the flow within the pump system. It is to be noted that, as used herein, variable frequency drives are to include adjustable frequency drives (AFDs), Variable Speed Controllers (VSCs) or something similar, which operate to control electronic motor speed.\nPump speed and pressure represent important pumping system parameters, in addition to flow, which can cause the pump to operate at less than its most efficient level. Even more disadvantageously, less than optimal operating parameters may cause the pump and motor to work harder and thus wear out quicker, thereby shortening the pump\"\"s operational lifetime. According, it is highly desirable to provide a computer-controlled variable frequency device (VFD) controller which utilizes computer algorithms and sensor inputs to control flow, speed, pressure and performance of a pumping system by monitoring motor, pump and system parameters and controlling pump output via speed variations. It is also advantageous to obtain a controller operative to identify and report pump or system anomalies to a technician, to facilitate investigation and correction of any abnormalities before any serious damage to the pumping unit occurs.\nA controller for controlling operating parameters associated with fluid flow, speed or pressure for a centrifugal pump for pumping fluid, wherein at least one sensor is coupled to the pump for generating a signal indicative of a sensed operating condition. The controller comprises a storage device for storing data indicative of at least one operating condition and a microprocessor in communication with the sensor and operative to perform an algorithm utilizing the at least one sensor signal and the stored data indicative of the at least one operating condition to generate a control signal, wherein the control signal is indicative of a correction factor to be applied to the pump.\nThere is also disclosed a method for automatically controlling operating parameters associated with a centrifugal pump according to an algorithm for pumping fluid to a discharge outlet, comprising the steps of storing in memory data values corresponding to predetermined operating conditions, obtaining sensor measurements indicative of current operating conditions, utilizing the sensor measurements and the stored data values to determine calculated data values corresponding to the current pump operating conditions, and comparing the calculated data values with the stored data values and generating a control signal indicative of a correction factor to be applied to the pump when the calculated data values differ from the stored data values by a predetermined amount."}
{"text": "Neurostimulation systems are devices that generate electrical pulses and deliver the pulses to nerve tissue to treat a variety of disorders. Spinal cord stimulation (SCS) is an example of neurostimulation in which electrical pulses are delivered to nerve tissue in the spine typically for the purpose of chronic pain control. Other examples include deep brain stimulation, cortical stimulation, cochlear nerve stimulation, peripheral nerve stimulation, vagal nerve stimulation, sacral nerve stimulation, etc. While a precise understanding of the interaction between the applied electrical energy and the nervous tissue is not fully appreciated, it is known that application of an electrical field to spinal nervous tissue can effectively mask certain types of pain transmitted from regions of the body associated with the stimulated nerve tissue. Specifically, applying electrical energy to the spinal cord associated with regions of the body afflicted with chronic pain can induce “paresthesia” (a subjective sensation of numbness or tingling) in the afflicted bodily regions. Thereby, paresthesia can effectively mask the transmission of non-acute pain sensations to the brain.\nNeurostimulation systems generally include a pulse generator and one or more leads. The pulse generator is typically implemented using a metallic housing that encloses circuitry for generating the electrical pulses, control circuitry, communication circuitry, a rechargeable battery, etc. The pulse generation circuitry is coupled to one or more stimulation leads through electrical connections provided in a “header” of the pulse generator. Specifically, feedthrough wires typically exit the metallic housing and enter into a header structure of a moldable material. Within the header structure, the feedthrough wires are electrically coupled to annular electrical connectors. The header structure holds the annular connectors in a fixed arrangement that corresponds to the arrangement of terminals on a stimulation lead.\nThe terminals of a stimulation lead are electrically coupled to the connectors within the header by manually pushing the proximal end of the stimulation lead into the header of the pulse generator. If the stimulation lead is improperly inserted into the header, the terminals of the stimulation lead will not be properly aligned with the connectors of the header and electrical pulses will not be properly conducted through the lead to electrodes for stimulation of tissue of the patient."}
{"text": "This application describes example systems and methods for placing onto a mammal's lower body, leg and thigh limbs, a set of paired bands to measure the lower body locomotion (for bipedalism, upright locomotion) and, if desired, additional arm strapped forearm and arm paired bands on the upper body (for complex motion, crawling, and in other applications, or four calf limb quadrupedalism locomotion). Each band contains multiple MEMs force sensors that measure muscle circumferential pressure at multiple positions, along with Earth's magnetic and gravitational fields. On-band data processing and networked, intra-band RF connectivity from multiple limbs, can be used to produce simple, energy-optimized, least-action metrics of mammal locomotion of two, interrelated functions defined as 1) Track, being a walking, running or other dynamic, from forward footpath creation over firm surfaces such as the ground, water-floats, or snow and ice, and 2) Balance, being required to efficiently move the lower body and/or upper body under Track. These technology metrics include assessment of locomotion related neurological functionality of body, limbs, and muscle disorders.\nDetails of Track differences in maintaining Balance are related to other physiological disorders, including detecting mental precursors, such as for fall-down behavior. An example system can be worn by many simultaneously and uniquely identified users, with metrics and location of each user being displayed on a laptop through additional RF connectivity, such as in recreational and professional sports. The system's on-band data processing units using battery supplied power, integrate the sensing to determine kinetic and potential energy of the body locomotion over time in a method that integrates out the aperiodic motion of the upper body and extended appendages, about center of mass (ACM), and uses the residual motion to measure the periodic center of mass (CM) locomotion from a known point.\nThe example systems and methods are self-calibrating, using calibration data collected when standing and jumping, at positions facing compass points. A GPS system can be incorporated for continuous motion measurement, to be used for calibration of the locomotion when GPS satellite data is available, and to establish the initiation geolocation point when beginning operation in GPS-denied regions for navigation, or to relatively geolocate multiple players in sport activities.\nThe body movement data, when combined with the band pressure and gravity data, removes the effects of the aperiodic, nonlinear locomotion, and leaves the residual movement for determining Track and Balance through a Euler-Lagrangian representation of the Equations of Motion.\nThe example systems and methods can be embodied for many applications (including the use of metric feedback for improved Balance and Track optimization), such as for human and non-human (e.g., horse) sports trainers, elder health care providers, and physical therapists for prevention and recovery from injuries of the back, knee, spine, and brain, such as from stroke, Alzheimer and Parkinson diseases, and dementia, as well as for other uses in mammal navigation and location purposes."}
{"text": "The invention relates to a fabric which exhibits a puffed design effect and method for the same. More particularly, the invention relates to the provision of a lightweight puffed embroidered design fabric and method by which a multitude of different designs may be created on fabrics by forming puffed design elements in a controlled manner according to prescribed embroidered patterns. The puffed embroidered fabric designs are relatively unlimited in their variety, and according to the fabric and method of the present invention, may be made on a commercial basis in large quantities.\nHeretofore, fabric designs have been created which are commonly referred to as three-dimensional, or puckered, fabric designs. The three-dimensional or puckered effect in the fabric design may be accomplished in a number of ways. For example, single ply fabric may be woven or knitted with areas of the fabric provided with yarns that are more shrinkable than the remaining yarns in desired areas. The shrinkable yarns are subjected to shrinkage by either hot washing, chemical, or other treatment, causing a swelled area in the fabric and a patterned effect such as in U.S. Pat. Nos. 2,738,566 and 3,071,165.\nU.S. Pat. No. 3,359,610 discloses a patterned puckered fabric achieved by weaving two plies of fabric together with an elastic yarn. The elastic yarn which joins the two plies of fabric together is shrunk creating puckering at the outer layer.\nPuckered or three-dimensional design fabrics have also been formed by simultaneously weaving two plies, and interweaving at spaced intervals, one of the plies, which may be shrunk, with the other ply which is not shrinkable. When the shrinkable ply shrinks, small wrinkles or puckers will form in the unshrinkable layer as taught by U.S. Pat. No. 2,231,388. Puckered two-ply fabrics have also been made in which puckering of the fabric is obtained by selective shrinking of certain of the threads of the woven two-ply fabric rather than shrinking one ply and not shrinking the other as described above. Such a fabric is shown in U.S. Pat. No. 2,401,829 wherein shrinkable warp or weft threads are interwoven between two fabric plies to form small puckers. However, the design possibilities in such fabrics are very limited due to the fact that the puckering pattern is established by interweaving two plies at selected points. Thus the designs are limited by the nature of the weaving process itself, and also by the fact that only certain types of designs with little or no variations may be created in this method. The creating of designs by weaving and interweaving two plies is a very time consuming and expensive process which is not commercially attractive. Moreover, the weight of the fabric produced by interweaving two plies results in fabric which has only very limited use and may not be used where a lightweight puckered fabric may be needed. The two-ply fabrics are typically dried on a tenter frame under tension which may take out shrinkage, particularly in the crosswise dimension of the fabric owing to the tensioning of the fabric between the tenter fames. Thus the puckered pattern is not easily varied or determined by the weaving or the drying of the two-ply fabric in the prior art methods and fabrics.\nU.S. Pat. No. 3,195,489 discloses a three-dimensional patterned fabric in which an outer layer is made from an elastic material and sewn onto a base layer. The two fabrics are then stretched and a design pattern is stitched to the stretched panels. Afterwards, the tension is released causing the stretched panel to contract and exhibit a pucker in the outer fabric layer. The ornamental fabric is mainly suitable for women's articles of clothing, namely swimsuits in which elastic panels are desired.\nAccordingly, an object of the present invention is to provide a method in which a wide variety of puffed designs may be created in a lightweight fabric.\nAnother object of the invention is to provide a puffed embroidered design fabric and method wherein a wide variety of embroidered designs can be created in the face of fabric which exhibit puffed design elements according to a prescribed well-controlled pattern.\nStill another object of the invention is to provide a puffed embroidered design fabric and method wherein a pattern of puffed elements is created in the embroidered design and the resultant fabric is lightweight and suitable for wearing apparel and other lightweight fabric applications.\nAnother object of the invention is to provide a puffed embroidered fabric and method which are well adapted to the production of a multitude of different designs on a commercial basis.\nStill another object of the invention is to provide a puffed embroidered design fabric and method wherein an embroidered pattern is stitched into a front fabric and a lightweight cotton gauze back fabric uniting the fabrics by interlocked embroidery stitches which control formation of puffed design elements in the embroidered pattern upon washing and drying of the fabric.\nYet another object of the invention is to provide a reversible fabric puckered and/or puffed on both faces in a predetermined pattern, as well as a method of making such a fabric, particularly on an embroidery machine."}
{"text": "Many products and services have only a limited time of usability. For example, if there are vacant seats on an airplane flight, when the flight takes off, the usability of (and the revenue from) these seats is lost. As another example, the usability of (and the revenue from) a hotel room is lost if the room remains unoccupied overnight. A third example would be a table at a popular restaurant. As a lunch hour or evening passes, the revenue and value from an unused table is lost. Many people might be interested in taking advantage of such opportunities at the last minute if the price is low enough, if the service is convenient enough, or if the service is particularly scarce, but they may not have the opportunity, the interest, or the time to search an auction site for such products and services.\nWhat is needed is a system and method that, based on a set of known preferences and historic transactions with customers, could cull from a customer database the customers most likely to close a transaction, and thus allow fleeting inventories to be sold off before they expire and allow customers to get access to extremely convenient, highly sought after, or well priced services."}
{"text": "The conventional method of conveying glass through a refractory wall within a glass melting tank is to incorporate the use of a channel made of refractory block. The channel or throat is generally located at the bottom of a bridgewall of the furnace, and is utilized to flow the molten glass from a melting chamber into an adjacent refining chamber, or for flowing the glass out of the furnace to be cooled in a forehearth. When there is molten glass on both sides of the wall dividing the various chambers, the throat or channel could be as simple as a hole in the wall. However, the action of the flowing molten glass tends to corrode the roof of the refractory channel, and typically the first part of the furnace to wear out is the throat area.\nThus, the effective life of glass melting furnaces is limited by the corrosion of the throats connecting the melting and fining zones or chambers, and those throats connecting the fining and refiner zones or chambers. The fining zone may be operated at temperatures above 1500.degree. C., so the walls and throats associated therewith suffer severe corrosion. The roof of the throats suffers the greatest corrosion, and even fused zirconia block is corroded away over time until the refractory is so thin that a glass leak can occur.\nRefractory metal pipes, such as molybdenum pipes, have been used in the past to convey some glasses from the bottom of a tank through a sidewall thereof into a distribution channel, such as shown in U.S. Pat. No. 4,029,887 to Spremulli. The flow of certain molten glasses through the moly pipe results in very small corrosion to the pipe per se. However, where the pipe passes through the sidewall, the corrosion of the surrounding refractory can be significant. The mechanism creating such high corrosion is known in the industry as upward drilling. That is, whenever there is a horizontal refractory surface with glass flowing therebeneath, bubbles in the glass will rise up until they hit the surface. The bubbles then tend to enhance corrosion due to a surface tension gradient on the bubbles, and the bubbles then in effect drill a hole into the refractory.\nThe dense corrosion products formed between the flowing glass and the corroded refractory then flow away, and fresh glass enters the area between the pipe and the refractory to repeat the corrosion process. When a molybdenum pipe is passed through a hole in the sidewall of a glass tank to flow molten glass from the tank, such as in the Spremulli patent, the refractory above the pipe will continue to corrode away even though the glass flows through the pipe. The corrosion of the refractory continues because there is nothing to prevent the dense corrosion products from flowing away. Eventually, there is very little refractory separating the glass from a water cooled outer rim of a refractory metal flange surrounding the pipe, and the process must be shut down before a glass leak occurs.\nU.S. Pat. No. 4,365,987 to Boettner discloses the use of water cooled molybdenum flanges to prevent leaks between adjacent portions of a glass delivery system. However, the Boettner patent is primarily directed to a refractory metal glass delivery system of molybdenum for controlling the flow of glass from a furnace to a forehearth by means of a flow control device incorporated therein. The inlet to the flow control system is at a level above the bottom of the furnace, and no means are provided for inhibiting or virtually eliminating the corrosion of the refractory block about the inlet end of the delivery tube. Accordingly, the corrosion of the refractory wall through which the pipe passes limits the life of the delivery system.\nOver the years attempts have been made to install throats within glass furnaces protected by refractory metal conduits such as molybdenum. However, such attempts have not been completely successful since no effort had been made to protect the refractory material surrounding the conduit from the corrosive effects of the glass on the outside of the conduit adjacent the surrounding refractory material. The use of a water-cooled flange about the conduit has helped to prevent complete flow through of the molten glass along the length of the conduit, but has not succeeded in preventing the corrosion of the refractory material adjacent the inlet and outlet ends and along the upper surface of the conduit. The present invention is directed to overcoming these deficiencies.\nIt thus has been an object of the present invention to provide method and apparatus for virtually eliminating or inhibiting the corrosion of refractory material about a refractory metal delivery system within a glass melting furnace."}
{"text": "Present day consumers typically own several electronic devices specifically designed for portability and on-the-go use, including, for example, a mobile phone or smart phone, a portable music player like an iPod® or an MP3 player, a tablet, a portable gaming unit, a camera, and the like. Each of these devices requires frequent recharging. Such electronic devices typically utilize a cable for connecting the device to a power source, such as a wall outlet, a car charger, an airplane charger, or a computer. However, a separate cable is usually required for each power source. Moreover, different electronic devices often utilize different connection ports and interfaces such that a single charging cable is not compatible with multiple devices. Accordingly, a tech-savvy consumer, with several electronic devices, will usually have multiple charging cables to keep track of. Even then, the consumer may be without sufficient power to recharge a phone due to bad weather or a power outage, or may not always be in a place where a power source is readily available, or even if so, may not have the appropriate cable or adapter available to use with a particular power source.\nWith traditional power sources, such as those noted above, it is often difficult to charge multiple devices at the same time, especially where each device requires a separate charging cable. For example, a car charger port may only handle a single cable at a time. Adaptor devices are available on the market for connecting multiple devices to a power source at the same time—for example, a two-to-one or three-to-one car charger splitter. However, such adapters are often only compatible with certain interfaces. Moreover, such adapters are separate from portable power sources and tend to be bulky.\nSimilarly, connection interface attachments are also available for adapting a charging cable for use with a variety of devices for recharging from a power source, each requiring a different interface connection. However, such attachments are usually separate small pieces, and therefore difficult to keep track of when not in use. Further, use of such attachments does not solve the problem presented by the need to charge multiple devices at the same time, from the same power source, as oftentimes, only one attachment can be used with a charging cable at a time.\nPortable power chargers exist that permit recharging of electronic devices when a standard power source is not readily available. For example, portable power chargers are illustrated and described in co-pending U.S. application Ser. No. 13/571,992, filed Aug. 10, 2012, and Ser. No. 13/682,985, filed Nov. 21, 2012, which share common inventors with the present application and which are incorporated herein by reference. Some existing power charger devices usually cannot charge multiple devices at the same time, either due to limited capacity or connectivity options. Even if multiple devices may be attached to the power charger at the same time, the charger may prioritize how the devices are recharged—i.e., it will charge one device first and then the second, and so on. However, this approach takes a long time to recharge all devices and risks not having sufficient charge remaining in the charger for fully charging the second device.\nFurther, some portable charger devices will not permit recharging from the charger when the charger is itself being recharged or connected to an external power source. Such devices require the charger unit to be disconnected from a power source before a charge will be passed on to a device connected to the charger, or require the charger unit to be fully charged first before any device connected to the charger unit can then be recharged.\nAdditionally, such portable charger devices typically require a dedicated input port for recharging the internal battery and a separate output port dedicated for recharging electrical devices from the internal battery. More particularly, such charging devices often require multiple output ports for recharging multiple electronic devices at the same time. The addition of extra charging ports compromises the size and design of the charger unit, for example, a unit with a dedicated input port and two or more output ports would need to be larger than a charger unit with just a single port due to the need to properly arrange the electronics for operation of the charger as desired.\nIn view of the foregoing, there is a need for a charger that can be used to charge a variety of electronic devices, including but not limited to smart phones, mobile phones, data tablets, music players, cameras, camcorders, gaming units, e-books, Bluetooth® headsets and earpieces, GPS devices, and the like, either individually or collectively in various combinations. Additionally, there is a need for such a charger that is portable, has a compact size, and is easy to use in various conditions and locations to charge one or more electronic devices simultaneously, or recharge the internal battery of the charger unit for future on-the-go use, including but not limited to in a house or office, a car or an airplane. Still further, there is a need for a portable charger having a port that can act both as an input port for recharging an internal battery unit in the charger and as an output port for recharging an electronic device connected to the charger. Still further, there is a need for a portable charger that can recharge the internal battery from an external power source (either from an AC power source or a DC power source) at the same time as an electronic device connected to the charger, even while both the external power source and the electronic device are connected to the charger through the same port. Still further, there is a need for a portable charger unit in a compact size that has increased functionality for a user requiring a portable source of power. Accordingly, it is a general object of the present invention to provide a portable charger that improves upon conventional power chargers currently on the market and that overcomes the problems and drawbacks associated with such prior art chargers."}
{"text": "The present invention relates to a control lever locking mechanism and more particularly relates to a locking mechanism which is activated by opening the door and deactivated by closing the door of the vehicle.\nControl levers which are convenient to the operator when seated are often times an obstacle to his entry and exit, and when bumped, may be moved to an actuated position initiating operation of a power control device connected to the vehicle thus creating a safety hazard."}
{"text": "The present invention relates to memory sells for integrated circuits, and, more particularly to memory cells whose retention elements integrate vertically stacked bipolar transistor transistors (“BJT”) with CMOS transistors.\nThe scale down tendency of the integrated circuit industry is causing an increasing technological barrier. One of the solutions is three-dimensional integrated circuits. These integrated circuits increase the manufacturing complexity significantly.\nFIG. 1A shows an ordinary prior art six-transistor memory cell. As can be seen, it has two complementary bit lines, a bit line and a negative bit line. These two complementary bit lines are also tied to differential sense amplifier that has the ability to speed up the loaded bit lines.\nMemory blocks hold 50-60% of dies area. Consequently, compacting the area of memory is equivalent to scaling down the size of the chip. The most serious effort to reduce the SRAM area (Static Random Access Memory) has been focused on the effort to reduce the number of transistors per memory cell. Such proposed solutions include the one-transistor SRAM (U.S. Pat. No. 6,765,830) and the Thyristor based SRAM (U.S. Pat. No. 6,944,051). All such proposed solutions, however, are not really SRAM, but rather are DRAM (Dynamic RAM) with SRAM interface. These memory cells require ongoing refresh cycles and significant amount of area for control and emulation of the SRAM interface. Consequently, these memory cells do not have the same applicability as an SRAM. Although these memory cells may achieve approximately 80% of area saving per cell, only less then 40% area saving will hold for full block of such memory with all the peripheral overhead.\nIt is noted that some of the SRAM cells designed in the past have demonstrated bipolar transistor cells or integration of such BJT (Bipolar Junction Transistors) with MOS transistors or other passive devices (U.S. Pat. No. 4,926,378 and U.S. Pat. No. 4,845,674).\nThere is a compelling need to have an apparatus that will significantly improve the compactness of integrated circuits, especially for SRAM. There is a further compelling need to have an apparatus that will also reduce the power requirements for such integrated circuits."}
{"text": "Infrared (IR) thermal cameras can be used in a number of different situations, for example, when inspecting or surveying complex electrical systems such as transformers, switchgears etc., or water carrying systems such as heat exchangers, radiators etc. IR cameras are used for capturing and storing thermal radiation data. This thermal radiation data may then be displayed and viewed as thermal images and analyzed in order to, for example, find faulty electrical wirings or couplings, leaking water pipes, etc.\nHowever, various procedures and methods are being used in order to properly analyse the thermal radiation data and/or the thermal images of the IR camera, and these are not necessarily particularly intuitive and easily understandable by a user of the IR camera. The analysis of the thermal radiation data and/or the thermal images of an IR camera may also be a time-consuming task and may thus preclude a user of an IR camera from making decisions, predictions, and/or recommendations to clients while being on site and performing the IR imaging."}
{"text": "Experimental and clinical observations have supported the concept that the hypothalamus plays a key role in the regulation of adenohypophysial corticotropic cells' secretory functions. Over 25 years ago it was demonstrated that factors present in the hypothalamus would increase the rate of ACTH secretion by the pituitary gland when incubated in vitro or maintained in an organ culture. However, a physiologic corticotropin releasing factor (CRF) was not characterized until ovine CRF (oCRF) was characterized in 1981. As disclosed in U.S. Pat. No. 4,415,558, the disclosure of which is incorporated herein by reference, oCRF was found to be a 41-residue amidated peptide. oCRF lowers blood pressure in mammals when injected peripherally and stimulates the secretion of ACTH and .beta.-endorphin.\nRat CRF (rCRF) was later isolated, purified and characterized; it was found to be a homologous, amidated hentetracontapeptide as described in U.S. Pat. No. 4,489,163, the disclosure of which is incorporated herein by reference. The formula of human CRF has now been determined to be the same as that of rCRF, and the terms rCRF and hCRF are used interchangeably.\nA CRF analog was subsequently developed having a high alpha-helical forming potential which is also a 41-residue amidated peptide. It is commonly referred to as AHC (alpha-helical CRF) and is described in U.S. Pat. No. 4,594,329, the disclosure of which is incorporated herein by reference. Other analogs containing D-isomers were developed, such as those shown in U.S. Pat. No. 5,278,146.\nSynthetic rCRF, oCRF and AHC stimulate ACTH and .beta.-endorphin-like activities (.beta.-END-Li) in vitro and in vivo and substantially lower blood pressure when injected peripherally. Antagonists of these compounds are disclosed in U.S. Pat. No. 4,605,642, issued Aug. 12, 1986, the disclosure of which is incorporated herein by reference. Cyclic CRF analogs exhibiting biopotency have been developed as disclosed in U.S. Pat. No. 5,245,009 (Sep. 14, 1993) and in pending U.S. patent application Ser. No. 78,558, filed Jun. 16, 1993, now U.S. Pat. No. 5,493,006.\nSince the foregoing discoveries, the search for improved CRF analogs has continued."}
{"text": "The present invention relates to rotary machines, and in particular, to an air deflector for an air cycle machine.\nAir cycle machines are used in environmental control systems in aircraft to condition air for delivery to an aircraft cabin. Conditioned air is air at a temperature, pressure, and humidity desirable for aircraft passenger comfort and safety. At or near ground level, the ambient air temperature and/or humidity is often sufficiently high that the air must be cooled as part of the conditioning process before being delivered to the aircraft cabin. At flight altitude, ambient air is often far cooler than desired, but at such a low pressure that it must be compressed to an acceptable pressure as part of the conditioning process. Compressing ambient air at flight altitude heats the resulting pressurized air sufficiently that it must be cooled, even if the ambient air temperature is very low. Thus, under most conditions, heat must be removed from air by the air cycle machine before the air is delivered to the aircraft cabin.\nTo condition the air as needed, air cycle machines include a fan section, a compressor section, and a turbine section that are all mounted on a common shaft. The compressor receives partially compressed air from the aircraft and further compresses the air. The compressed air then moves through a heat exchanger and is cooled by the fan section. The air then moves through the turbine section where it is expanded for use in the aircraft, for example, for use as cabin air. The turbine section also extracts energy from the air and uses the energy to drive the fan section and the compressor section via the common shaft.\nAir cycle machines also include bearings that are positioned around the common shaft. The bearings are cooled by passing a cooling air flow through a cavity that is adjacent the bearing. The cooling air flow then exits the cavity and is discharged from the air cycle machine into an ambient. The cooling air flow is limited in that it can only cool the bearing using convective heat transfer. The cooling air flow is further limited in that the cooling air flow in the cavity flows through a center of the cavity, meaning a majority of the cooling air flow does not flow across a surface of the bearing."}
{"text": "Internal combustion engines for heavy-duty and industrial applications are typically fuelled by diesel. However, the use of natural gas as an alternative to diesel is of increasing interest. Natural gas is relatively abundant and relatively cheap, and can, in principle, provide similar levels of power to diesel whilst producing lower particulate and nitrogen oxide (NOx) emissions.\nNatural gas can be used in place of diesel to fuel a compression-ignition engine, in which combustion of the fuel occurs as a result of compression of the air-fuel mixture in the cylinder. However, because natural gas has a higher auto-ignition temperature than diesel, it can be necessary to initiate combustion with a pilot injection of diesel fuel before introducing the natural gas to the combustion chamber.\nIn one type of natural gas-powered engine, known as a high-pressure direct injection (HPDI) engine, both natural gas and diesel are injected directly into the combustion chamber. Due to the space constraints in an engine cylinder head, it is desirable to inject both fuels using one fuel injector per cylinder. This requires a fuel injector that is specially adapted to keep the two fuels separate within the injector, and to deliver independently the respective fuel at the appropriate time.\nOne such ‘dual fuel’ injector is described in International Patent Application Publication No. WO 00/15956. In this example, a fuel injector with a concentric twin nozzle arrangement is provided. Inner and outer valve needles are engageable at their lower ends with respective valve seats to control the flow of fuel through respective inner and outer sets of outlets. The outer valve needle controls the injection of natural gas through the outer set of outlets, and the inner valve needle controls the injection of diesel through the inner set of outlets. The outer valve needle is tubular to accommodate the inner valve needle, and the inner set of outlets is formed at a tip of the outer valve needle.\nThe inner and outer valve needles are controlled independently by two electromagnetic control valves, which are configured to control the pressure of a control fluid (normally diesel fuel) within respective control chambers for the inner and outer valve needles. The control chambers receive the upper ends of the respective needles, so that changing the pressure of the control fluid in each control chamber changes the downward (closing) force on the corresponding needle. Gas or diesel fuel pressure acts on downwardly-facing thrust surfaces of the respective needles to generate an upward (opening) force on the needle. When the pressure of the control fluid in a control chamber is relatively high, the downward force is greater than the upward force and the respective needle remains seated, and when the pressure of the control fluid is relatively low, the upward force overcomes the downward force and the respective needle opens to permit fuel injection through the respective set of outlets.\nEach control chamber is connected to a source of control fluid at relatively high pressure. Each control valve is operable to connect the respective control chamber to a low-pressure drain for the control fluid. In this way, opening of each control valve causes a reduction in the pressure of the control fluid in the corresponding control chamber, resulting in opening of the corresponding valve needle.\nThe injector is mounted in a bore in the cylinder head of the engine. A main body of the injector extends through the bore, so that the tip of the injector protrudes into the respective combustion chamber. The maximum diameter of the cylinder head bore, and hence the diameter of the main body of the injector, is constrained by the limited space available in the cylinder head. The electromagnetic actuators used to actuate the control valves are generally too large to be accommodated in the main body of the injector. Instead, the actuators and control valves are typically mounted in a side-by-side configuration at the top of the fuel injector, so that they protrude above the cylinder head.\nInjectors of this type therefore have a relatively bulky top end, which takes up space in the region above the cylinder head. Furthermore, the injectors must be supplied with two fuels from two independent fuel supply systems. Most conveniently, both fuel supply systems are mounted on top of the cylinder head. However, because some of the space above the cylinder head is reserved for the valve train of the engine and other engine components, the space available for the fuel rails, supply pipes and other components of the fuel supply systems is limited. It would therefore be desirable to provide fuel injectors and fuel injection systems in which more compact and/or more convenient fuel supply arrangements are possible."}
{"text": "1. Field of the Invention\nThis invention relates to an improved support for a lithographic plate, and more particularly to a support for a lithographic plate comprising a plate of aluminum or an alloy thereof (hereafter referred to as aluminum) which is characterized by the depth of a grain defined by the center line averge roughness (Ra) as well as the average diameter of pits in the grain and the distribution thereof.\n2. Description of the Prior Art\nHeretofore, an aluminum plate has been widely used as the support of a photosensitive lithographic plate. The surface of the aluminum plate is made coarse to give better adhesion to a photosensitive layer to be disposed on the plate or to provide higher wettability (water retention) with wettening water during printing. This process of surface roughening is called graining and the rough surface obtained is a grain. Two conventional methods of graining are mechanical graining and electrochemical graining. Ball graining which is one example of the former method is a very old method which has found wide utility in small-scale graining operation. However, this method is not efficient since it does not permit continuous operation. An industrially applicable mechanical graining that is substituted for ball graining is brush graining which rotates a brush or roller having steel wires or synthetic resin hairs, optionally in the presence of a sand or abrasive, on an aluminum plate. Whether ball graining or brush graining is used, a grain having various degrees of surface roughness and shape can be produced by properly controlling the kind and shape of the sand, mesh, graining period and brush movement. However, as reported in TAGA (Technical Association of Graphic Arts) Proceedings, pp. 262-276, 1972 mechanical graining leaves residual sand on the surface of the grain as part of the grain. In addition, as shown in the picture taken by a scanning electron microscope in tha technical paper, mechanical graining of an aluminum plate provides a worked surface having a very complicated three-dimensional configuration. Therefore, if a lithographic plate is produced from a presensitized plate having a photosensitive layer disposed on a mechanically grained aluminum plate, the sensitive material penetrates too deeply into the pits in the grain to be removed easily, thus often providing a stained non-image area.\nMany studies have been made on electrochemical graining. British Pat. No. 831,998 as well as U.S. Pat. No. 3,072,546 and 3,073,765 disclose an aluminum support for a lithographic plate which is grained by A. C. electrolysis using hydrochloric acid as electrolyte. British Pat. No. 1,224,226 discloses a method of A. C. electrolysis of an aluminum plate in hydrochloric acid, followed by chemical etching and anodization. Japanese Patent Publication No. 27481/71 discloses a method of A. C. electrolysis of an aluminum plate in hydrochloric acid, followed by anodization.\nIt is generally known that A. C. electrolysis of an aluminum plate in an electrolyte mainly consisting of hydrochloric acid or nitric acid can provide the aluminum plate with a grained surface. The grain provided by electrochemical graining which is formed of grown pits has a crater-like or honeycomb stucture and is characterized by having straight and open pits as compared with the grain produced by mechanical graining described above. Another feature of electrochemical graining is that it provides a plate having deeper pits and a coarser grain then mechanical graining. The configuration and coarseness of the grain can be controlled by selecting the electrolyte and electrolytic conditions used. West German Patent (OLS) No. 2,650,762 describes the grain provided by electrochemical graining using hydrochloric acid or nitric acid as an electrolyte. A coarse-grained surface obtained by using nitric acid or an electrolyte mainly consisting of nitric acid produces pitting having a dual structure which comprises a pit provided by electrochemical etching plus an extremely small pit formed in its surface, but the opening is generally shallow. On the other hand, the use of hydrochloric acid or an electrolyte mainly consisting of hydrochloric acid provides pitting which is generally deep but the surface of individual pits is relatively smooth without a complex configuration as achieved by use of a nitric acid based electrolyte.\nAlthough an aluminum plate grained electrochemically has a by far a coarser surface than a mechanically grained aluminum plate, it has not yet been considered a preferred support having advantages over and replacing the conventional supports for lithographic plates. Among the problems yet to be solved with an electrochemically grained aluminum plate are its low affinity for ink, short running life (number of sheets of paper which can be printed from one plate), its low printing sensitivity and low developing speed."}
{"text": "The latest one chip MCU frequently has a RAM monitor function for debugging and tuning an internal program of the MCU in a state implemented on a chip (hereinafter referred to as MCU 100 with a built-in RAM monitor or simply referred to as MCU 100). FIG. 1 shows the configuration of a conventional MCU 100.\nThe RAM monitor function can be realized when the MCU 100 receives an external clock and external data and outputs data by way of an external interface terminal group 102.\nThe conventional MCU 100 comprises an external interface terminal group 102, a data transmit/receive section 104, a direct memory access control circuit (hereinafter referred to as DMA controller) 110, and also comprises a CPU 120, a ROM 130, a RAM 140, a peripheral module 1 (150) and a peripheral module 2 (160) respectively connected with one another by way of an internal bus formed of a data bus (DBUS) 106 and an address bus (ABUS) 108.\nThe DMA controller 110 comprises a GO flag set F/F 111, an AND gate 115, an address register (AReg) 112, a data write register (WReg) 113, a data read register (RReg) 114, and a timing controller 116.\nThe GO flag set F/F 111 is set upon completion of reception of an address/data from the transmit/receive section 104 and reset upon completion of a DMA transfer operation. An output signal from the GO flag set F/F 111 is inputted to the AND gate 115 together with a bus state recognition signal 101 and generates a DMA start signal 124. The bus state recognition signal 101 is a signal which takes an insignificant value (e.g. “0”) when the CPU 120 occupies the bus and the DMA is placed in a bus use prohibition state, and which takes a significant value (e.g. “1”) when the CPU 120 releases the bus and the DMA is placed in a bus use permission state.\nFurther, the timing controller 116 executes write/read relative to the registers, the RAM and the like by the DMA operation at a given timing at the time when a GO flag is set by the GO flag set F/F 111 and the bus state recognition signal 101 takes the significant value and the operation of the timing controller 116 is reserved until a starting condition is prepared.\nThe RAM monitor function is a function to read or rewrite contents of the RAM 140 and registers of the peripheral modules 1 and 2 inside the MCU 100 during the normal operation of the MCU 100, namely, during the operation of the MCU 100 by an application program (hereinafter referred to as AP) which is written in the ROM 130.\nIn the case of reading data, in FIG. 1, when an address of a RAM or addresses of registers, namely, the RAM 140, the peripheral module 1 (150) and the peripheral module 2 (160) which are provided in the MCU 100 and subjected to debugging are inputted from an external data input/output terminal (hereinafter referred to simply as external data i/o terminal) in synchronization with an external clock from the external interface terminal group 102, necessary data is set in the DMA controller 110 by way of the data transmit/receive section 104. The DMA controller 110 reads data in the RAM or registers serving as objects of set addresses by way of the internal bus, and outputs the read data to the outside by way of the data transmit/receive section 104 and the external interface terminal group 102.\nIn the case of writing data, when an address of a RAM or addresses of registers, namely, the RAM 140, the peripheral module 1 (150) and the peripheral module 2 (160) which are provided in the MCU 100 and subjected to debugging and data to be written are inputted from the external data i/o terminal in synchronization with the external clock from the external interface terminal group 102, necessary data is set in the DMA controller 110 by way of the data transmit/receive section 104. The DMA controller 110 rewrites data stored in the RAM or registers serving as objects of set addresses by way of the internal bus, and outputs a signal indicating the completion of execution of writing to the outside by way of the data transmit/receive section 104 and the external interface terminal group 102.\nHowever, in the conventional method, since a point (a state of program counter (PC 121) inside the CPU 120) where the data in the RAM and registers inside the MCU 100 are read or data is written on the RAM and registers is a time upon completion of setting the DMA controller 110 after setting data from the external interface terminal group 102, such a point can not be set from the outside, causing a problem of lack of flexibility of debugging and another problem of low accuracy of debugging.\nUnder the circumstances, an on-chip debugging method of a MCU which is improved in flexibility and accuracy of debugging has been desired."}
{"text": "It has been over 40 years since the discovery of the “Australia antigen” in hepatitis B virus (HBV) from hemophiliac patients in 1965 by Blumberg. However, a thorough understanding of the life cycle of HBV has remained elusive so far, and there is no powerful drug for the treatment of chronic hepatitis B. The life cycle of HBV is pretty complex. The basic steps involves coupling of virus outer coat membrane with specific receptor(s) on the cell surface, entry of the virus into the cell through fusion of the viral membrane with cell membrane or phagocytosis, disassembly of the viral outer capsid in inclusion bodies, nuclear entry of core protein together with HBV DNA, which is repaired into cccDNA, and then transcription of various viral mRNAs and translation of the viral mRNAs into viral proteins for the assembly and secretion of viral particles (see Seeger C, Mason W S. Hepatitis B virus biology. Microbiol. Mol. Biol. Rev. 2000; 64:51-68). The detailed studies of HBV life cycle depends on a good HBV infection model system. Several cell lines, e.g., HepG2 and Huh-7, have been established that can support HBV replication after transfection with HBV DNA plasmids. However, transfection using HBV plasmids introduces the HBV DNA directly into these cells and does not involve the early steps of the HBV life cycle, such as virus interactions with specific cell membrane receptors, virus uncoating, and transport of the HBV DNA into the nucleus. Thus, these cell lines have limited use in the studies of the HBV life cycle. So far, only a few cells are known that can support HBV infection and replication. These cells include primary human liver cells (including embryonic liver cells), primary liver cells from tree shrew, and the recently established HepRG cells. Due to scarcity of the available sources, stringent requirements for isolation and culture conditions, and short survival time once isolated (generally less than 1 month), human primary liver cells and tree shrew primary liver cells are not suitable for scientific studies that require long-term investigation or large-scale cell culture (Guha C, Mohan S, Roy-Chowdhury N, et al. Cell culture and animal models of viral hepatitis, Part I: Hepatitis B. Lab. Anim. (NY) 2004; 33:37-46). HepRG cell line was established in 2002. This cell line is a liver progenitor cell line, which can be infected by HBV after pre-differentiation induction using DMSO. This cell line is not commercially available, but is only used in the lab where it was initially established. HepRG cells require pre-differentiation induction before they can be used in HBV infection studies. Stability and reproducibility of this model, however, are yet to be verified (see Gripon P, Rumin S, Urban S, et al. Infection of a human hepatoma cell line by the hepatitis B virus. Proc. Nat. Acad. Sci. USA. 2002; 99:1565-15660). Therefore, there is no reproducible cell line that can be used for the study of the complete HBV infection cycle. This lack of a suitable model has greatly impeded in-depth studies of various steps in the life cycle of HBV, especially the studies of the viral specific cellular receptor and the HBV DNA nuclear entry mechanism in the early stages of the HBV life cycle. The lack of a suitable cell line for HBV infection has also hindered the efforts to screen for and develop novel anti-HBV drugs. Therefore, identification of cell lines that are naturally susceptible to HBV infection remains important for the advancement of HBV research and for the development of therapeutics for chronic hepatitis."}
{"text": "Modeling methods are known generally in the state of the art, for example from the following articles:\nHOSKING J. M. R. et al. “A statistical perspective on data mining” FUTURE GENERATION COMPUTER SYSTEMS, November 1997, ELSEVIER, Vol. 13, No. 2-3, pages 117-134, which describes the use of statistical methods for processing large volumes of data;\nROSSIUS R. et al. “A short note about the application of polynomial kernels with fractional degree in support vector learning” MACHINE LEARNING ECML-98. 10th European Conference on Machine Learning. Proceedings CHEMNITZ Apr. 21-23, 1998, pages 143-148, pertaining to the application of classification and forecasting methods.\nThe learning problem can be considered as a problem of finding dependencies using a limited number of observations. Thus, it is a question of choosing from a given set of functions f(x,α), α∈A, where A is a set of parameters, the one which best approximates the output.\nIf L(y,f(x,α) is a measure of the deviation between the real output y and the output predicted by the model f(x,α), it is thus necessary to minimize the effective risk:R(α)=∫L(y,f(x,α))dF(x,y)  (Eq. 1)while knowing that the joint probability distribution function F(x,y) is unknown and that the only available information is contained in the k observations (x1, y1) . . . ,(xk, yk) from the learning set.\nClassically, one determines the function that minimizes the empirical risk calculated on the basis of the learning set:                     Remp        =                              ∑                          i              =              1                        k                    ⁢                      L            ⁡                          (                                                y                  i                                ,                                  f                  ⁡                                      (                                                                  x                        i                                            ,                      α                                        )                                                              )                                                          (                  Eq          .                                           ⁢          2                )            \nOne then postulates that this function would be the best approximation of the function that minimizes the effective risk given by (Eq. 1).\nThe problem posed is to know the extent to which a system constructed on the empirical risk minimization inductive principle (Eq. 2) is generalizable, i.e., enables minimizing the effective risk (Eq. 1) including data that have not been learned.\nMathematically, a problem is said to be well-posed when it allows a unique solution and this solution is stable, i.e., that a small deviation in the initial conditions can only modify in an infinitesimal manner the form of the solutions. Problems that do not satisfy these conditions are referred to as ill-posed problems.\nIt occurs frequently that the problem of finding f satisfying the equality A,f=u is ill-posed: even if there exists a unique solution to this equation, a small deviation of the right-hand side of this equation can cause large deviations in the solution.\nAnd thus if the right-hand member is not exact (uε instead of u with ∥u−uε∥≦ε), the functions that minimize the empirical risk R(ƒ)=∥Aƒ−u∈∥2 are not necessarily good approximations of the desired solution, even if ∈ tends to 0.\nAn improvement in solving such problems consists in minimizing another so-called regularized functional of the form:R(ƒ)=R(ƒ)+λ(∈)Ω(ƒ)  (Eq. 3)where:                Ω(ƒ) is some functional belonging to a special type of operators referred to as regularizing;        λ(∈) is an appropriately chosen constant depending on the level of noise existing on the data.        \nOne then obtains a sequence of solutions that converge to the desired one as ∈ tends to 0. Minimizing the regularized risk rather than the empirical risk allows obtaining from a limited number of observations a solution that is generalizable to any case.\nIntroduction of the regularizing term makes it feasible to provide with certainty a unique solution to an ill-posed problem. This solution can be slightly less accurate than the classic one, but it possesses the fundamental property of being stable, thus endowing the results with greater robustness.\nThe methods for solving ill-posed problems demonstrate that there exist other inductive principles that enable obtaining a better regularization capacity than the principle consisting in minimizing the error made on the learning set.\nTherefore, the main objective of theoretical analysis is to find the principles making it feasible to control the generalization capacity of learning systems and to construct algorithms that implement these principles.\nVapnik's theory is the tool that establishes the necessary and sufficient conditions for a learning process based on the empirical risk minimization principle to be generalizable, leading to a new inductive principle referred to as the structural risk minimization principle. It can be demonstrated that the effective risk satisfies an inequality of the form:R(α)<Remp(α)+F(h,k)  (Eq. 4)where:                h is the Vapnik-Chervonenkis dimension of the set of functions f(x,α) among which the solution is sought;        k is the number of observations available for constructing the model;        F is an increasing function of h and a decreasing function of k.        \nIt can be seen immediately that, since the number k of available observations is finite, the fact of minimizing the empirical error is not sufficient for minimizing the effective error. The general idea of the structural risk minimization principle is to take into account the two terms of the right-hand member of the equation (Eq. 4), rather than only the empirical risk. This implies constraining the structure of the set of the functions f(x, α) among which the solution is sought so as to limit or even control the parameter h.\nAccording to this principle, the development of new algorithms enabling control of the robustness of learning processes would be advantageous."}
{"text": "1. Field of the Invention\nThe present invention relates to a light extraction substrate and an organic light-emitting device having the same, and more particularly, to a light extraction substrate which can realize a superior light extraction efficiency when applied to an organic light-emitting device, and an organic light-emitting device having the same.\n2. Description of Related Art\nIn general, an organic light-emitting diode (OLED) includes an anode, a light-emitting layer and a cathode. When a voltage is applied between the anode and the cathode, holes are injected from the anode into a hole injection layer and then migrate from the hole injection layer through a hole transport layer to the organic light-emitting layer, and electrons are injected from the cathode into an electron injection layer and then migrate from the electron injection layer through an electron transport layer to the light-emitting layer. Holes and electrons that are injected into the light-emitting layer recombine with each other in the light-emitting layer, thereby generating excitons. When the excitons transit from an excited state to a ground state, light is emitted.\nOrganic light-emitting displays including an OLED are divided into a passive matrix type and an active matrix type depending on the mechanism that drives the N*M number of pixels which are arranged in the shape of a matrix.\nIn an active matrix type, a pixel electrode which defines a light-emitting area and a unit pixel driving circuit which applies a current or voltage to the pixel electrode are positioned in a unit pixel area. The unit pixel driving circuit has at least two thin-film transistors (TFTs) and one capacitor. Due to this configuration, the unit pixel driving circuit can supply a constant current irrespective of the number of pixels, thereby realizing uniform luminance. The active matrix type organic light-emitting display consumes little power, and thus can be advantageously applied to high definition displays and large displays.\nHowever, as shown in FIG. 5, only about 20% of light generated from an OLED is emitted to the outside and about 80% of the light is lost by a waveguide effect originating from the difference in the refractive index between a glass substrate 10 and an organic light-emitting layer 30 which includes an anode 20, a hole injection layer, a hole transport layer, a light-emitting layer, an electron transport layer and an electron injection layer and by a total internal reflection originating from the difference in the refractive index between the glass substrate 10 and the air. Specifically, the refractive index of the internal organic light-emitting layer 30 ranges from 1.7 to 1.8, whereas the refractive index of indium tin oxide (ITO) which is generally used for the anode 20 ranges from 1.8 to 1.9. Since the two layers have a very small thickness ranging from 200 to 400 nm and the refractive index of glass used for the glass substrate 10 is about 1.5, a planar waveguide is thereby formed inside the organic light-emitting device. It is calculated that the ratio of the light lost in the internal waveguide mode due to the above-described reason is about 45%. In addition, since the refractive index of the glass substrate 10 is about 1.5 and the refractive index of the ambient air is 1.0, when the light is directed outward from the inside of the glass substrate 10, a ray of the light having an angle of incidence greater than a critical angle is totally reflected and is trapped inside the glass substrate 10. Since the ratio of the trapped light is up to about 35%, only about 20% of the generated light is emitted to the outside.\nIn addition, as shown in FIG. 6, in order to overcome the foregoing problem, in the related art, a low index grid (LIP) 50 is formed on the ITO anode 20. The grid 50 converts the direction of the light that travels in the waveguide mode to the front surface, thereby improving light extraction efficiency.\nFIG. 7 shows simulation results on the organic light-emitting device shown in FIG. 6. The effect of improving the light extraction efficiency is increased when the refractive index of the grid 50 is lower. However, there are problems in that almost no materials have a refractive index of 1.2 or less and that the price of a material is more expensive when the refractive index is lower. In addition, when the grid 50 is formed on the ITO anode 20, as shown in FIG. 6, a stepped portion is formed. Consequently, a leakage current may occur. In addition, the organic light-emitting device shown in FIG. 6 has the problem of difficult processing. For example, in some cases, the surface of the anode 20 which adjoins the organic light-emitting layer 30 is metamorphosed in the process of forming the grid 50 on the ITO anode 20, thereby changing the work function. Furthermore, holes are not injected into the organic light-emitting layer 30 through the portion of the anode 20 on which the grid 50 is formed, and the size of the electric field applied thereto is different from the surroundings, thereby decreasing the uniformity of the light generated.\nIn addition, as shown in FIG. 8, in the related art, a convex-concave structure 60 is disposed under the anode 20 (with respect to the paper surface), i.e. in the interface between the anode 20 and the glass substrate 10, in order to improve light extraction efficiency.\nAs described above, the anode 20 and the organic light-emitting layer 30 generally serve as one light waveguide between the cathode 40 and the glass substrate 10. Accordingly, in the state in which the anode 20 and the organic light-emitting layer 30 act in a waveguide mode, when the nanoscale convex-concave structure 60 which causes light scattering is formed on the surface that adjoins the anode 20, the waveguide mode is disturbed, so that the quantity of light that is extracted to the outside is increased. However, when the convex-concave structure 60 is formed below the anode 20, the shape of the anode 20 resembles the shape of the convex-concave structure 60 below the anode 20, thereby increasing the possibility that a sharp portion may be localized. Since the OLED has a stacked structure of very thin films, when the anode 20 has a sharp protruding portion, the current is concentrated in that portion, which is a cause of a large leakage current or decreases power efficiency. Accordingly, in order to prevent such deterioration of the electrical characteristics, a flat film 70 is necessarily added when the convex-concave structure 60 is formed below the anode 20. The flat film 70 is required to be made of a material that anode 20. When the refractive index of the flat film 70 is low, most of the light is reflected at the interface between the anode 20 and the flat film 70 and then is trapped between the anode 20 and the organic light-emitting layer 30, which is referred to as the waveguide mode. In addition, the flat film 70 is required to be as thin as possible. If the flat film 70 is too thick, more light may be unnecessarily absorbed, and the distance between the convex-concave structure 60 and the organic light-emitting layer 30 may become too large, thereby reducing the scattering effect.\nThe information disclosed in the Background of the Invention section is provided only for enhancement of (or better) understanding of the background of the invention, and should not be taken as an acknowledgment or any form of suggestion that this information forms a prior art that would already be known to a person skilled in the art."}
{"text": "It is known that many processes performed during the manufacture of a semiconductor device require heating of the workpiece. It is most desirable that this heating be uniform across the surface of the workpiece since this will tend to result in a uniform process being performed.\nIn the prior art, different heating configurations have been proposed to attain relative temperature uniformity. During a typical process, there is a temperature ramping phase where the temperature of the wafer is brought up rapidly to a predetermined level, and a steady state phase where the temperature is maintained at about this level. The temperature non-uniformity has tended to be greater during the ramping phase.\nOne approach to attaining a uniform temperature is called \"zone control\". In zone control, a plurality of radiant heating units such as infrared lamps are divided into zones and the lamps in a given zone receive the same power, while lamps in different zones may receive different power, to result in uniform temperature across the wafer. While zone control is advantageous in some applications, it requires complex electronic controls of the power applied to the lamps.\nA high temperature, rapid thermal process using flood heating is known. With flood heating, a temperature inversion may be observed wherein the edge of the workpiece is hotter during the ramping phase, while the center of the wafer is hotter during the steady state phase. It has been recognized that this is undesirable in the high temperature environment (e.g. 1150.degree. C.) of rapid thermal processing because the edge to center temperature differences at such process temperature create stresses which can damage the workpiece."}
{"text": "It is known in the art to manufacture urethanes by reacting hydroxyl group-containing organic compounds, such as alcohols or phenols, with carbon monoxide and certain nitrogenous compounds, using metal carbonyls, certain metal complexes, and certain metallic compounds, as catalysts. Thus, British Pat. No. 1,080,094, to Ibbotson, teaches the above-described process using as a catalyst a mixture of two or more metal halides, at least one of the said metal halides being a halide of a transition metal; British Pat. No. 1,092,157, to Gamlen et al., teaches using as a catalyst a metal complex compound containing at least one transition metal and at least one unsaturated hydrocarbon ligand in which the unsaturated system forms a bond with the metal; British Pat. No. 1,087,896, to Ibbotson, teaches using as a catalyst a mixture of one or more metals capable of existence in two or more valency states and one or more halides of metals other than the metal or metals first mentioned and also capable of existence in two or more valency states; U.S. Pat. No. 3,338,956. to Mountfield, teaches using as catalyst, metal carbonyls. Applicant has discovered that the above reaction can produce very high yields of diurethane with relatively little undesirable side products if it is carried out in two stages, the first stage utilizing one temperature and pressure, and the second stage being carried out at a higher temperature and pressure."}
{"text": "In the UMTS (Universal Mobile Telecommunications System) network, the specifications of long term evolution (LTE) have been drafted for the purpose of further increasing high speed data rates, providing lower delays and so on (see non-patent literature 1). In LTE, as multiple-access schemes, a scheme that is based on OFDMA (Orthogonal Frequency Division Multiple Access) is used in downlink channels (downlink), and a scheme that is based on SC-FDMA (Single Carrier Frequency Division Multiple Access) is used in uplink channels (uplink). Also, successor systems of LTE (also referred to as, for example, “LTE-advanced” or “LTE enhancement” (hereinafter referred to as “LTE-A”)) have been developed for the purpose of achieving further broadbandization and increased speed beyond LTE, and the specifications thereof have been drafted (Rel. 10/11).\nIn relationship to LTE-A systems, a HetNet (Heterogeneous Network), in which small cells (for example, pico cells, femto cells and so on), each having local a coverage area of a radius of approximately several tens of meters, are formed within a macro cell having a wide coverage area of a radius of approximately several kilometers, is under study. Also, in relationship to HetNets, a study is in progress to use carriers of different frequency bands between macro cells and small cells, in addition to carriers of the same frequency band.\nFurthermore, for future radio communication systems (Rel. 12 and later versions), a study is in progress to introduce a new mechanism for small cell discovery. To be more specific, a user terminal detects the detection/measurement signal that is transmitted from a small cell, and establishes synchronization, conducts measurements and so on. Note that this detection/measurement signal may be referred to as the “discovery reference signal” (DRS). When a channel state measurement signal (for example, the CSI-RS (Channel State Information Reference Signal) is included in the DRS, the user terminal detects and measures this channel state measurement signal based on the timing of synchronization."}
{"text": "The present invention relates to an educational device and includes a plurality of movable plates that simulate a shoe to provide not only play value, but is also educational in purpose in teaching the lacing and tying of shoelaces.\nEducational devices that also have play value are commonly available for use and instruction of children, and these prior known devices have taken various forms, depending upon the purpose in the teaching aspect thereof. Prior to the instant invention, some efforts have been made to provide educational devices that incorporate the configuration of a shoe, wherein the simulated shoe also included a provision for teaching the lacing and tying of shoelaces on the shoe, thereby enabling a child to conveniently practice the technique of shoelace tying at a relatively early age. Such an educational device as known heretofore is illustrated by way of example in U.S. Pat. No. Des. 142,670; and, as shown in this patent, a book comprising front and back covers and a plurality of pages are formed in a design that simulates a shoe. However, the play value of this prior known educational device and others similar thereto had limited interest for the child and was constructed in such a way as to be easily abused by a child in the use thereof.\nAlthough the prior known educational devices did have some educational and play value and in some instances accomplished the purpose intended, they were usually not constructed or formed in a way to retain the interest of a child; and, since such articles did not impart a realistic effect to the article being portrayed, the child usually quickly lost interest after a short period of use and play thereof."}
{"text": "There is the need in deer hunting and the like to place one's self above the ground in order to command a superior field of vision of the hunting area and to keep out the deer's line of sight. Therefore, tree stands are a preferred method of hunting for deer. Building an ordinary stand takes a considerable amount of time and effort. Once the stand is built it is more or less permanent. Although suitable for the instant purpose, it soon becomes impractical. The building materials are subject to deterioration; the deer population in the immediate area may decrease necessitating a relocation of the stand; also, a permanent stand might attract other hunters to one's favorite spot.\nThere has accordingly long been a need to develop a prefabricated tree stand which may quickly and efficiently be placed in and removed from the tree. While some chair stands have been developed, there has remained the need for a compact, integrated unit which features easy manipulation, lightweight construction, portability and versatility in other uses."}
{"text": "1. Field of the Invention\nThis invention relates to a peptide which has an affinity for gp120, HIV (human immunodeficiency virus) envelope protein.\n2. Description of the Related Art\nThe therapy for HIV infection is usually chemotherapy, such as the nucleotide derivative AZT (3′-azido-3′-deoxythmidine). This AZT therapy or protease inhibitor, which was later developed, prolongs the life of HIV patients, but there are some problems, these are derived from the chemotherapy itself.\nThe problems are shown as follows: The first is chronic poisoning due to long term administration, the second is the appearance of an HIV virus resistant to the medicine during the therapy, the third is the appearance of malignant tumors in prolongation of the HIV patient's life, the fourth is that the recovery of the immune system can not be monitored, the fifth is that there is not a method to monitor treatment effect, etc. Since such chemotherapy is not basic therapy for HIV infections, most people anticipate the development of a vaccine.\nGenerally, the vaccine is an inactive treatment (in active vaccine) of a microbe of viruses; a weak activity virus which loses pathogenesis or a pseudo virus (live vaccine) which has no fatal effects to humans. However, although HIV itself is natively a weak activity virus, it is well known to stay long after having once invaded the body. In addition, the host cell of HIV is mainly a lymphocyte, which controls the immune system; furthermore, HIV spreads over hemophilic patients through blood-preparation. From these finding, even if it is assumed that we selected either an in-active or a weak vaccine, the development of an HIV vaccine has many problems with safety.\nAccordingly, an HIV vaccine is being developed which utilizes a part of the HIV envelope protein and inhibits further infection.\nFrom such an idea, many researchers performed an epitope analysis of gp120 in the HIV enveloped protein, and then, watched the V3 region (3rd hypervariable region) of gp120 as an epitope. But it was a true hypervariable region [Palker T. J., et al., Proc. Natl. Acad. Sci. USA 85:2709-2713, 1988; Rusche J. R., et al., ibid 85:3198-3202, 1988; Gouddsmit J., et all., ibid 85:4478-4482.1988; Matsushita S., et al., J. Virol. 62:2107-2114, 1988]. After this, a vaccine which used a part of this region as antigen was administrated to an HIV infected monkey as an infection inhibitory experiment, but the effectiveness has not yet been reported.\nAs well as this, Tam et al. devised further antigenecity for the above-mentioned peptide antigen (Tam et. al., Japanese patent publication (Tokuhyo) No. H 3-503539), but have not yet had success because in most parts of the V region, particularly in the V3 which is a convenient region for antibody preparation, mutation or deletion occurs.\nIn addition, a neutralized antibody, which inhibits the infection against lymphocyte, is developed. For example, in Japanese Patent Application No. 63-171385, after the production of a monoclonal antibody by using a part of the above mentioned peptide as antigen, a method is reported, which produces anti HIV chimera antibodies on hybridization due to genetic engineering at the level of the protein as the Fab′ itself. But, although with such neutralizing antibodies it is possible to inhibit HIV infection to the lymphocyte at laboratory level, an antibody that can be used practically has not yet been developed.\nAs mentioned above, chemotherapy has some problems; drug tolerance in the virus and side effects in the host, another idea to solve the problem of removing the virus from the body is by plasmapheresis. Although this method to remove the HIV virus by using a pore size membrane filter (smaller than virus size) for plasmapheresis has been definitely proposed it is not yet possible to make a uniform pore size membrane. It is also possible that the pores will become clogged during plasmapheresis resulting in the deterioration of the membrane due to pressure. As mentioned above there are many technical problems which have to be settled. So, a method to use CD4 derived from human lymphocyte having specific affinity to HIV, as absorbed carrier in column for plasmapheresis is also proposed. It cannot be used as a medical procedure because of the lost affinity due to decay by autoclave treatment. In addition, there are also methods using thermostable molecules, a high molecule polymer or color ligand as an affinity carrier to HIV. However, as these molecules do not originally have specific binding ability to HIV, they cannot be used because they bind to blood ingredients faster than to HIV.\nIn this way, aiming at the development of an HIV treatment medicine, research to produce a vaccine and neutralizing antibody is flourishing, but useful medicine has not yet been developed.\nThe inventors paid attention to this present situation and developed a superior peptide to have the same degree or more affinity for gp120 compared to antibodies and to be resistant to autoclave treatment, and have already made a patent application (Japanese Patent Application No. H 8-351474 and Japanese Patent Application No. 8-351475). This peptide basically consists of a three amino acid sequence, but from a study of the sequel, we found that an affinity to gp120 of this peptide deteriorated by number and a kind of the amino acid which ranged in it. So, we knew that we needed to develop a more stable peptide."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a frequency modulator, and in particular, to a digital frequency modulator for accumulating input digital signals and changing the phase of the accumulated signal.\n2. Description of the Related Art\nIn general, a wireless communication system transmits on a radio channel a relatively high-frequency carrier with a relatively low-frequency input signal. This is called modulation. In particular, when the information of the input signal is contained in the frequency and phase of the carrier signal, the modulation is called frequency modulation (FM).\nAn analog frequency modulator illustrated in FIG. 1 is usually used as a frequency modulator. The analog frequency modulator is comprised of a gain controller 110, a digital-to-analog converter (DAC) 120, an FM filter 130, a voltage controlled oscillator (VCO) 150, and a frequency synthesizer 140.\nReferring to FIG. 1, the gain controller 110 multiplies an input digital signal Vin by a predetermined gain G. The DAC 120 converts the output of the gain controller 110 to an analog signal. The FM filter 130 filters the analog signal and feeds the filtered signal as a voltage control signal to the VCO 150. Meanwhile, the frequency synthesizer 140 generates a carrier signal with a predetermined carrier frequency. The VCO 150 controls the frequency of the carrier signal received from the frequency synthesizer 140 with the voltage control signal, thereby outputting an FM signal.\nAs described above, the analog frequency modulator effects FM by converting an input digital signal to an analog signal and then controlling the VCO 150 with the analog signal.\nA distinctive shortcoming of the typical analog frequency modulator is that the VCO being an analog device makes performance optimization difficult due to performance variation and use of additional analog devices for stabilizing the characteristics of the VCO increases cost. The VCO has a non-linear frequency output characteristic with respect to an input signal and exhibits very different characteristics depending on devices. Thus to achieve desirable FM characteristics, compensation or adjustment is required for the VCO. As a solution to this problem, implementation of a digital frequency modulator has been considered. However, the digital frequency modulator also has the limitations of difficulty in achieving desirable resolution with respect to various frequency inputs and implementation complexity."}
{"text": "The use of lightweight composite materials in vehicles, such as aircraft, provide an improved strength-to-weight ratio that translates to fuel savings and lower operational costs. Composite materials, however, do not readily conduct away extreme electrical currents and electromagnetic forces generated by lightning strikes. Furthermore, composite or metallic structures used on vehicles are typically assembled using metallic fastener systems, which are conductive and therefore create electromagnetic effect (EME) design considerations. For example, fastener sparking modes must be designed for lightning conditions that include hot particle ejection and arcing between the fastener and surrounding structures."}
{"text": "This invention relates to a convenient wall-mounted holder for shaving razors, specifically for holding a dispenser of multiple razors. These razors can be of the type having replaceable blades, but more typically entire razors are sold, each of which is disposable.\nAt present, a package of multiple razors is sold in a plastic bag. A user must remove one razor at a time and remove a protective blade guard. Since the razor can be used more than one time before becoming dull, the current razor is placed haphazardly around the sink or ledge in a shower. Since the package of new razors is usually at some distant location from the sink or shower, it becomes problematic how to obtain a new razor, how to deal with the current razor, and what to do with a depleted razor.\nPrior art shown in U.S. Pat. No. 6,598,609 to Gibbs discloses a housing for combining a can of shaving cream with a shaver and replacement shaver heads, but this housing must be placed on a shelf. Prior art shown in U.S. Pat. No. 6,796,455 to Schmidt discloses a wall-mounted blade dispenser whose blades are not of the type used for shaving. Prior art shown in U.S. Pat. No. 6,869,053 to Adams discloses a suction cup with a ring which could be used for hanging a razor."}
{"text": "Conventional projection systems typically use 3 primary colours (red, green and blue) for the reproduction of colour images. The colour gamut that can be produced by an additive combination of the 3 primary colours nevertheless is limited. The colour gamut depends on the dominant wavelength (hue) and on the excitation purity (saturation) of each of the primary colours. The visible colours, the primary colours and the gamut of produced colours are usually represented in a chromaticity diagram e.g. CIE 1931 chromaticity diagram or CIE 1976 U.C.S. (Uniform Chromaticity Scale).\nOne solution to produce a wider colour gamut is by increasing the excitation purity of the primary colours, or in other words to narrow the spectral pass band of each of the three primary colours. An alternative solution to produce a wider colour gamut may be the use of more primary colours, such as e.g. using 4, 5, 6 or more colours instead of 3 colours. In printing devices, where subtractive mixing of colours is performed (by different ink cartridges), the commonly used colours are cyan, magenta, yellow and black (CMYK). In special printers, dedicated for certain types of printing, additional colours such as e.g. indigo also are used. In display technology, several techniques are known to apply more than three primary colours. One option is to generate a wider colour gamut by generating more than three primary colours by filtering them from a white illumination source and modulating them sequentially according to image data using a single chip modulation system. In such a system, typically the amount of illumination from the illumination source that is not used for displaying the image is relatively high resulting in a less efficient system. In another option, the wider colour gamut is generated by guiding more than three primary colours each to a primary colour dedicated modulator, where the primary coloured sub-beam is modulated. The latter typically results in complex and expensive systems. An example of such a system is a display system using four primary colours, each primary colour sub-beam modulated by its own primary colour dedicated modulator, as described for example in “Four primary colour projection display” by Roth and Caldwell in SID 05 Digest (2005) 1818.\nAdding more primary colours also introduces an electronic puzzle as to how to describe the colour space, as each colour to be displayed needs to be produced by combining the different primary colours, i.e. by subtracting in the case of printing, or by adding in the case of a display system. A method of handling the data conversion for generating images to be displayed using more than three primary colours has been described e.g. in European Patent EP 0 897 641 B1 by BARCO N. V."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for roasting skewered food which can quickly cook a large amount of skewered roasted food such as skewered roasted chickens, skewered roasted giblets and the like prepared by heating pieces of meat, fish, vegetables and the like set on a skewer.\n2. Description of Related Art\nThere has been known the method of cooking skewered roasted food by setting pieces of meat, fish, vegetables and the like on a skewer and heating the materials, but for cooking skewered roasted food such as skewered roasted chickens, skewered roasted giblets and the like which can be provided and sold at pubs, restaurants, or other types of stores, persons hoping to engage in the job are required to work as apprentices at skilled persons for several years. In addition, because scales of the pubs, restaurants, and other stores have been increasingly larger, it is required to cook a large amount of skewered roasted food quickly within a short period of time.\nTo satisfy the need as described above, there has been proposed the apparatus for roasting skewered food which can quickly cook a large amount of skewered roasted food within a short period of time without requiring the long time job training, in which a number of skewer holders are mounted on an endless rotary chain with a prescribed space in-between, raw materials of skewered food prepared by setting food materials on a skewer are held on the skewer holders in order, and the raw materials of skewered food are carried along heating devices for heating the raw materials of skewered food to prepare skewered roasted food (Refer to, for instance, Japanese Patent Laid.-Open Publication No. 2004-57463).\nIn the conventional type of apparatus for roasting skewered food as described above, when a raw material of skewered food is held on skewer holder, it is required to hold the raw material of skewered food in the substantially horizontal posture. However, it is difficult for an operator to hold the raw material of skewered food in the horizontal posture, and because it is required to hold the raw material of skewered food with the tip portion of the skewer facing against the operator, the food materials are touched by the operator's fingers, which is not desirable from the sanitary point of view.\nFurthermore, there is no step of immersing a raw material of skewered food in sauce during the process of carrying the raw material of skewered food along heating devices with an endless rotary chain for heating, and therefore, after the raw material of skewered food is heated to roast, it is necessary to immerse the raw material of skewered food in sauce and heat the raw material of skewered food again, so that the work load is rather heavy and the period of time required for cooking skewered roasted food is disadvantageously long.\nIn addition, because linking pins or the like are used for mounting the skewer holders on the endless rotary chain, when the number of skewer holders is to be changed or any skewer holder is exchanged with a new one because of a failure or for some other reasons, it is required to remove the endless rotary chain from sprockets and then dismount skewer holders, and therefore, the work for dismounting skewer holders is laborious and requires a long period of time."}
{"text": "1. Field of the Invention\nThe present invention is directed to a therapy apparatus for treating a patient with focused acoustic waves.\n2. Description of the Prior Art\nTherapy systems which employ acoustic waves are utilized, for example, for treating stone complaints (lithotripsy), tumors and bone pathologies (osteorestoration). These therapy systems usually are in the form of completely equipped work stations which, in addition to comprising a source of focused acoustic waves with application means for introducing the generated acoustic waves into a patient to be treated, include a locating system and a patient supporting device, whereby the source and the patient supporting device are adjustable relative to one another. During the treatment, the region of the patient to be treated is first localized with the locating system, and the region to be treated is then positioned by adjusting the patient supporting device with the patient thereon and the source relative to one another so that this region is situated in the focus of the acoustic waves. The region to be treated is then charged with acoustic waves in the required way. Although therapy systems having a locating system functioning exclusively on the basis of ultrasound are suitable and are frequently used, therapy systems are normally preferred that have an x-ray locating system, since it is desirable in nearly all applications (and indispensable in many applications) to be able to locate a region to be treated with x-radiation. Some therapy systems have an x-ray locating system, in addition to an ultrasound locating system, since the latter is capable of supplying additional information.\nTherapy systems having an x-ray locating system are extremely expensive. This does not represent a disadvantage when the clinic in which the therapy system is operated has an adequately large number of patients in order to enable an economical operation of the therapy system. In the interests of the patient, however, it would be desirable to be able to have a high-performance therapy system of the type described above in smaller clinics or in a doctors's office as well. In the case of a therapy system equipped with an x-ray locating device, however, the comparatively high price thereof does not allow an economical utilization of the therapy system under those conditions."}
{"text": "1. Field of the Invention\nThe present invention relates to improved molding compositions. More particularly, the present invention relates to improved molding compositions comprising an intimate blend of a high nitrile polymer and a formaldehyde compound. These molding compositions are especially adapted for fabrication into molded packaging material which have a very low level of extractable HCN.\n2. Description of the Prior Art\nIn recent years it has been discovered that certain polymeric nitrile resins are especially suitable for packaging applications because of their excellent water and oxygen barrier properties. Such polymers are described at length in U.S. Pat. Nos. 3,451,528, 3,615,710, 3,426,102 and British Pat. No. 1,186,361, among others. Even more recently it has been discovered that although such nitrile polymers have excellent barrier properties they may be unsuitable for certain packaging applications because they contain trace, though detectable, amounts of hydrogen cyanide (HCN) which may be extracted by and impart a taste to the contents of the package. The amount of HCN in such nitrile polymers will vary with the nitrile monomer, the total nitrile content of the polymer, the polymerization method used to prepare the polymer, the processing and thermal history of the polymer and other factors. Moreover, it has been found that the amount of extractable HCN in nitrile polymers may be increased after the polymer has been subjected to forming operations wherein the polymer has been heated in order to soften and/or melt the polymer prior to forming it into shaped articles.\nThe trace amount of extractable HCN present in such nitrile packaging materials is very low and presents no known health or safety problems. In fact, the amount of extractable HCN in such nitrile polymers is lower than that found in many foods. In this regard it should be noted that HCN is a natural component in many foods and many other foods such as cereals, cocoa, ham, bacon and sausage, which are fumigated with HCN (prussic acid), are permitted to contain from 50 to 200 parts per million of HCN. (See the Food Additive Regulations of the FDA (Page 31, Subpart D, Paragraph 121,1072) as published in the Feveral Register: Dec. 23, 1965; 30 F.R. .15912 and the 1962 Public Health Service Publication 956 \"Drinking Water Standards\" ).\nHowever, as stated above, the HCN in certain nitrile packaging materials may, in certain instances, be extracted and impart a taste to the contents of the package. The problem of HCN extraction as it affects taste, is of concern in the packaging of beverages which are in prolonged intimate contact with the container. In these packaging applications the probability of HCN extraction is very high.\nThus, there exists in the art a need for improved nitrile polymer molding compositions which can be fabricated into packaging materials which are substantially free of extractable HCN. This need is fulfilled by the present invention which provides a molding composition adapted for the fabrication of packaging materials with a very low extractable HCN content. More particularly, molding compositions of the present invention are especially suitable for the preparation of packaging materials such as bottles, cans, jars, etc., which are used in the packaging of beverages, including carbonated beverages, where the probability of HCN extraction is very high."}
{"text": "Mobility management and power management are two essential functions in a wireless system. For a wireless device to be reached at any time or location, there is a need for the wireless device to support a low rate connection while not actively engaged in a voice call or data session. Mobility management involves messaging between the wireless device and the radio access network (RAN) related to connection management functions, for example, location, paging, cell reselection, and handoff.\nEffective power management is essential to conserve battery power. Power saving modes in a wireless device, such as idle mode or sleep mode, decrease the wireless device functions to a minimum so as to maintain some connectivity to the network while conserving power. Each radio access technology (RAT) standard stipulates its own mobility and power management functions. To perform mobility and power management efficiently, wireless standards have introduced state machines at the network layer and/or at the Media Access Control (MAC) layer.\nModern wireless devices, such as smartphones, netbooks, tablets can operate in heterogeneous network regions comprising network nodes supporting different RATs, e.g., GSM/EDGE/3G base stations, WiFi Access Points). Most wireless devices are capable of multi-mode operation (e.g., GSM/EDGE/3G/WiFi/Bluetooth/NFC), i.e., there are multiple radios in a wireless device capable of simultaneously accessing different Radio Access Networks (RANs) with each RAN supporting one or more RATs.\nFuture wireless device configurations will include even more modes of radio operation including future technologies such as LTE-Advanced as well as sensor radios and cognitive radios. In a multi-mode wireless device, radios that are not actively transmitting (e.g., in IDLE mode) use up some bandwidth and power for state management functions. The battery power usage increases with more radios in a wireless device, even if they are in their lowest power operating state. Further, placing radios in their lowest power operating state results in potential delays when the wireless device has to respond to an incoming voice call or data session. Therefore, there is a need for efficient low power operation modes for a multi-radio wireless device, such as the systems and methods discussed below in the present disclosure."}
{"text": "Unless otherwise indicated herein, the materials described in this section are not prior art to the claims and are not admitted to be prior art by inclusion in this section.\nA typical cellular wireless network includes a number of base stations that radiate to define wireless coverage areas, such as cells and cell sectors, in which wireless communication devices (WCDs) (also known as user equipment devices (UEs)), such as cell phones, tablet computers, tracking devices, embedded wireless modules, and other wirelessly equipped communication devices, can operate. In turn, each base station may be coupled with network infrastructure that provides connectivity with one or more transport networks, such as the public switched telephone network (PSTN) and/or the Internet for instance. With this arrangement, a WCD within coverage of the network may engage in air interface communication with a base station and may thereby communicate via the base station with various remote network entities or with other WCDs served by the base station.\nIn general, a cellular wireless network may operate in accordance with a particular air interface protocol or “radio access technology,” with communications from the base stations to WCDs defining a downlink or forward link and communications from the WCDs to the base stations defining an uplink or reverse link. Examples of existing air interface protocols include, without limitation, Orthogonal Frequency Division Multiple Access (OFDMA (e.g., Long Term Evolution (LTE)), Code Division Multiple Access (CDMA) (e.g., 1×RTT and 1×EV-DO), Wireless Interoperability for Microwave Access (WiMAX), and Global System for Mobile Communications (GSM), among others. Each protocol may define its own procedures for registration of WCDs, initiation of communications, handoff between coverage areas, and functions related to air interface communication.\nIn accordance with the air interface protocol, each coverage area may operate on one or more carrier frequencies or “carriers.” More particularly, the base station that radiates to define a given coverage area may support one or more frequency bands, such as the 800 MHz band (one or more frequency ranges around 800 MHz), the 1.9 GHz band (one or more frequency ranges around 1.9 GHz), and the 2.5 GHz band (one or more frequency ranges around 2.5 GHz), and may provide service on one or more carrier frequencies with each supported band. In a frequency division duplex (FDD) arrangement, different carrier frequencies are used for the downlink than the uplink. Whereas, in a time division duplex (TDD) arrangement, the same carrier frequency is used for the downlink and uplink and is allocated over time among downlink and uplink communications.\nOn each carrier frequency in a coverage area, the coverage area may also define a number of air interface channels for carrying information between the base station and the WCDs. These channels may be defined in various ways, such as through frequency division multiplexing, time division multiplexing, and/or code-division multiplexing, for instance. By way of example, each coverage area may define a pilot channel, reference channel or other resource on which the base station may broadcast a pilot signal, reference signal, or the like that WCDs may detect as an indication of coverage and may measure to evaluate coverage strength. As another example, each coverage area may define an uplink control channel or other resource on which WCDs may transmit control messages such as registration requests and access requests to the base station. And each coverage area may define a downlink control channel or other resource on which the base station may transmit control messages such as system information messages and page messages to WCDs. Each coverage area may then define one or more traffic channels or other resources for carrying communication traffic such as voice data and other data between the base station and WCDs.\nWhen a WCD first powers on or enters into coverage of the network, the WCD may scan for and identify a strongest pilot or reference signal and may register with the network by transmitting a registration request or attach request to a base station providing that signal. This registration process may serve to notify the network of the WCD's presence in a particular coverage area and to facilitate network authentication of the WCD. Once registered, the WCD may then operate in an idle mode in which the WCD monitors a downlink control channel to receive overhead information and to check for any page messages. In the idle mode, the WCD may have no assigned traffic channel resources on which to engage in bearer communication.\nWhen the network has a communication (such as a voice call or other traffic) to provide to a WCD that is registered with the network but is operating in the idle mode, the network may page the WCD in an effort to then facilitate assigning traffic channel resources to the WCD. In particular, the network may transmit on the downlink a page message addressed to the WCD. Assuming the WCD receives this page message, the WCD may then transmit to the network a page response message on the uplink. And upon receipt of the page response message, the network may then assign traffic channel resources to the WCD, for use to carry the communication, thus transitioning the WCD to a connected or active mode in which the WCD can engage in the communication.\nLikewise, when an idle WCD seeks to initiate a communication (such as to place a voice call or engage in other bearer communication), the WCD may transmit on the uplink to the base station an origination or connection request, and the network may then assign traffic channel resources to the WCD for use to carry the communication, similarly transitioning the WCD to a connected or active mode in which the WCD can engage in the communication."}
{"text": "Internal combustion engines, including diesel engines, gasoline engines, gaseous fuel-powered engines, and other engines known in the art exhaust a complex mixture of air pollutants. These air pollutants are composed of gaseous compounds such as, for example, the oxides of nitrogen (NOX). Due to increased awareness of the environment, exhaust emission standards have become more stringent, and the amount of NOX emitted from an engine may be regulated depending on the type of engine, size of engine, and/or class of engine. In order to ensure compliance with the regulation of these compounds, some engine manufacturers have implemented a strategy called Selective Catalytic Reduction (SCR).\nSCR is a process where gaseous or liquid reductant (most commonly a solution of urea solid and water) is added to the exhaust gas stream of an engine and is adsorbed onto a catalyst. The reductant reacts with NOX in the exhaust gas to form H2O and N2, which can be safely released to the atmosphere. Although SCR can be an effective method for reducing NOX, it can be difficult to ensure that enough reductant has been injected to adequately reduce the amount of NOX present within the exhaust gas stream, without unnecessarily wasting reductant and releasing unused reductant or byproducts thereof to the atmosphere.\nTo help control the amount of NOX emitted to the atmosphere, most exhaust systems are equipped with a sensor that detects the level of NOX in the tailpipe of an engine and generates a corresponding signal. Based on this signal, the amount of reductant injected into an exhaust flow can then be trimmed to precisely reduce a desired amount of NOX.\nAlthough relatively successful at helping to control the amount of NOX emitted to the atmosphere, most NOX sensors currently available are cross-sensitive to the reductant used to reduce NOX. That is, NOX sensors on the market today generate signals indicative of the presence of both NOX and reductant. As such, when either NOX or reductant levels in an exhaust flow go up, the signal value from the corresponding NOX sensor also goes up, and vice versa. Thus, it can be difficult to determine which constituent is currently being detected by the sensor, and over-injecting of reductant can occur in some situations when the injections are unknowingly based on a detected high level of reductant instead of NOX.\nOne attempt to help ensure that reductant injections are correctly based on detected levels of NOX is described in WO Publication Number 2006/000877 of Crane et al. that published on Jan. 5, 2006 (the '877 publication). The '877 publication discloses an exhaust system that utilizes a nonselective NOX/NH3 sensor in conjunction with feedforward and feedback control. The sensor is periodically interrogated by introducing a positive pulse in an ammonia feed rate. A positive response from the sensor to the positive pulse indicates ammonia slip is occurring (i.e., that the sensor is currently detecting NH3). A negative response from the sensor indicates NOX breakthrough (i.e., that the sensor is currently detecting NOX). Based on the positive responses, a controller inhibits feedforward control for a time and reduces reductant injections. Based on the negative response, the controller adjusts feedforward control to increase reductant injections. Feedback control is utilized between sensor interrogations to trim feedforward commands for injection based on sensor output under normal conditions.\nAlthough perhaps somewhat effective at determining if a sensor is detecting NOX or NH3, the exhaust system of the '877 publication may be unresponsive and wasteful. That is, the exhaust system of the '877 publication requires dedicated time for the sensor to be interrogated, this time causing delays in normal operations of the system. In addition, the exhaust system of the '877 publication requires extra pulses of reductant that are used only for sensor interrogation (i.e., the pulses do not help to reduce NOX), which can be expensive and decrease a time between servicing because of the additional exposure to reductant.\nThe exhaust system of the present disclosure solves one or more of the problems set forth above and/or other problems of the prior art."}
{"text": "With the repaid development of wireless communication systems, the radiating antennas with low cost and dual-band operation are in demand for various applications, such as wireless local area network (WLAN). The WLAN systems unitize the unlicensed bands, the ISM band at 2.4 GHz and 5.2 GHz only for the industrial, scientific and medical applications, to reduce the cost of the network building. The slot antenna has been investigated since the 1940s and many research studies have discussed the feeding structure, the bandwidth characteristics and the radiation phenomenon. The slot is etched on the ground plane of the substrate and radiation from both sides of the substrate is achieved by a microstrip or CPW-fed in a resonator cavity. In order to derive the maximum beam at the broadside direction, the length of the slot antenna should be limited to one wavelength. The technology of applying two different slot loops in the aperture was proposed to create dual bands. FIG. 12 shows a conventional slot antenna comprising a plurality of etched slots 62 in the ground plate 61 of the substrate 60 and a sub-antenna 64 extends from the ground plate 61 and to the opening 63 of the slot 62 so as to be the signal feedline. Another conventional slot antenna is disclosed in FIG. 13 and the slot 72 is etched as a T-shaped slot in the ground plate 71 of the substrate 70. A T-shaped feedline 73 is located in the inner periphery of the slot 72. Nevertheless, none of the conventional slot antennas is satisfied.\nThe present invention intends to provide the CPW-fed slot antenna to perform the dual-band characteristic and miniaturize the slot size utilizing the E-like feeding structure and adding the matching stub at the slot edge and four floating patches on the backside of the substrate. According to the measured results, the bandwidth of the lower resonant frequency distributes from 2.38 GHz to 2.60 GHz and one of higher resonant frequency ranges from 5.13 GHz to 5.73 GHz. The above techniques are effective and realizable for the design of the CPW-fed slot antenna."}
{"text": "The present invention relates to the packaging of integrated circuit dies and, more particularly, to heat dissipation in lead frames.\nVery similar processes are utilized to manufacture a large variety of integrated circuit packages or chips. A starting substrate, e.g., a thin wafer of silicon or gallium arsenide, is masked, etched and doped during several process steps, with the type, number and order of the steps depending on the type of integrated circuit being manufactured, to form a number of dies or separate integrated circuits thereon. The dies are singulated or separated with a wafer saw and then packaged individually to form integrated circuit packages or chips.\nEach die typically has a back that is devoid of circuitry and a front having integrated circuitry formed thereon. The integrated circuitry is accessible via die wire bonding pads that may be arranged in a variety of configurations on the face or edges of the die.\nDuring the packaging process, each die is attached to a corresponding lead frame. The lead frames are processed in groups, typically in strips or in a continuous coil form. Each strip contains multiple, e.g., six, lead frames and is several, e.g., nine, inches long. A coil contains a larger number of lead frames because it is a much longer continuous strip of material. Each lead frame strip or coil typically is a metal frame designed to support several dies for packaging and to provide the leads for the final integrated circuit package.\nA typical lead frame strip or coil is produced from metal sheet stock, e.g., a copper or nickel alloy, by subjecting the sheet stock to a stamping or photochemical etch process to form lead fingers, and side rails having pilot holes therein. The side rails and pilot holes therein facilitate transport and indexing of the lead frame strip or coil by automated packaging machinery. For example, a lead frame strip, individually, or a lead frame coil, as a whole, can be removably clipped to carrier belts that move the lead frame strip or coil through the various manufacturing stages.\nEach lead frame has a plurality of lead fingers for connection to the die bonding pads. The lead fingers typically are spot plated with palladium, gold or silver. The conductive leads are plated to provide a metallic surface to which wires may be bonded, as a bond wire usually will not stick directly to lead frame material, such as copper or nickel or alloys thereof.\nAs the demand for thinner packages increases, it becomes desirable to make the die-lead frame assembly thinner, which can be accomplished by reducing the thickness of the lead frame. However, thinner lead frames exacerbate the problems of heat dissipation both during manufacture and use. Particularly, thinner lead frames will, in the absence of other modifications, be less capable of dissipating the tremendous amount of heat that might build up either during manufacture or during operation.\nProblems of heat dissipation have also been encountered in view of increased power demands in the context of integrated circuit packages or chips. Particularly, as integrated circuit packages and chips are increasingly used in operating environments that notably increase the power provided thereto, the need for efficient and effective heat dissipation becomes more and more apparent.\nTo attend to the task of heat dissipation, manufacturers often provide supplementary media for dissipating heat away from the lead fingers. These supplementary media have often taken the form of small metal elements that are connected to portions of a lead frame.\nConventionally, this attachment is often achieved by using small rivets or tappets to interconnect a portion of a lead frame with the heat dissipating element. However, many disadvantages have been encountered with such an arrangement. Particularly, the use of separate interconnecting members, for the purpose of attaching a heat dissipating element to the lead frame, results in extra components that require undesirable additional expenditures of materials and cost.\nIt is also known to impart a significant thickness to the heat dissipating element in order to increase heat dissipation capabilities. However, by imparting a significant thickness to a heat dissipating element, limitations are subsequently placed on the overall thickness of the integrated circuit package or chip being manufactured.\nGenerally, the ends of lead fingers, when juxtaposed adjacent to one another, form the outline of a square or other rectilinear shape which might or might not be occupied by additional material. One type of conventional heat dissipating arrangement involves the placement of a quantity of thermally conductive material within this rectilinear space, along with thermally conductive connections to the siderails or other conducting portions of the lead frame. However, since the areal extent of such a heat dissipation arrangement is limited by the dimensions of the rectilinear space, it is often the case that only a significantly small heat dissipating element is realizable, with the result that heat dissipating capabilities may not be as great as desired.\nThe present invention broadly contemplates, in accordance with at least one presently preferred embodiment, a lead frame apparatus comprising: a lead frame; and an arrangement for dissipating heat generated at the lead frame, the heat dissipating arrangement being uninterruptedly connected to the lead frame.\nIn accordance with another presently preferred embodiment, the present invention broadly contemplates lead frame heat dissipating apparatus comprising at least one element for dissipating heat generated at a lead frame, the at least one element comprising an arrangement for directly and uninterruptedly connecting with a lead frame.\nIn accordance with another presently preferred embodiment, the present invention broadly contemplates lead frame apparatus comprising an arrangement for directly and uninterruptedly accommodating at least one external element for dissipating heat generated at the lead frame.\nIn accordance with yet another presently preferred embodiment, the present invention broadly contemplates a method of making lead frame apparatus, the method comprising the steps of: providing a lead frame; providing an arrangement for dissipating heat generated at the lead frame; and uninterruptedly connecting the heat dissipating arrangement to the lead frame.\nFinally, in accordance with yet another presently preferred embodiment, the present invention broadly contemplates lead frame apparatus comprising: a lead frame; an arrangement for dissipating heat generated at the lead frame; and an arrangement for mechanically connecting the lead frame with the heat dissipating arrangement; the connecting arrangement consisting essentially of constituent portions of the lead frame and the heat dissipating apparatus."}
{"text": "The present invention relates to a displacement sensor and, in particular, relates to such a device which senses the displacement or vibration of a living body. The present invention is applicable to a stethoscope for medical purposes, and mechanocardiogram.\nConventionally, an acoustic stethoscope has been used for diagnosing a living body. A prior stethoscope places a diaphragm on a surface of a living body, and the vibration of the diaphragm is listened to by a doctor through a stethoscope tube.\nTherefore, a prior stethoscope has the disadvantage that only a single doctor can listen to the sound of vibration of heart valves of a living body, but it can not be listened to by a plurality of persons. Further, a prior stethoscope has the disadvantage that a doctor cannot listen to the sound of a heart beat due to a low frequency less than 20 Hz and the output of the stethoscope can not be recorded.\nOn the other hand, a prior electrocardiogram records pulsation or drive electrical signals of a heart, but it does not record an actual mechanical movement of a heart. Therefore, if there is something wrong with a heart such as valve desease and arteriosclerosis, it could not be detected by a prior electrocardiogram. Conventional mechanocardiograph sensors for detection of heart movement are difficult to treat due to their bad treatability."}
{"text": "F152 (LL-Z1640-2) (1) is a zearalenone-like macrolide, isolated from shake flask fermentation, crude extracts of which inhibited the ciliated protozoan Tetrahymena pyriformis (see, McGahren et al. J. Org. Chem. 1978, 43, 2339). It was reported that initial biological studies using this natural product failed to yield any particularly interesting activities.\n\nAfter initial isolation and reporting of this compound, several other groups explored the possibility of preparing additional derivatives and/or further exploring their biological activity. For example, scientists at Merck reported that F152 and certain isomers thereof inhibit the phosphorylating enzyme Map/Erk kinase (MEK) and thus are useful for the treatment of certain cancers and other diseases characterized by the formation of neoangiogenesis (see, GB 323 845). Other groups have also reported derivatives of F152 having activity as tyrosine kinase inhibitors, which are useful, for example, for the treatment of cancer and inflammatory disorders (see, EP 606 044; WO 00/38674; JP 8-40893; WO 96/13259; U.S. Pat. Nos. 5,728,726; 5,674,892; 5,795,910). Each of these groups, however, was only able to obtain F152 and derivatives thereof by fermentation techniques and by modifications to the natural product, respectively, and thus were limited in the number and types of derivatives that could be prepared and evaluated for biological activity. Additionally, although F152 and certain derivatives thereof have demonstrated potent in vitro activities, these compounds are biologically unstable (for example, they are susceptible to enone isomerization in mouse and human plasma), thereby limiting the development of these compounds as therapeutics for the treatment of humans or other animals.\nClearly, there remains a need to develop synthetic methodologies to access and examine the therapeutic effect of a variety of novel analogues of F152, particularly those that are inaccessible by making modifications to the natural product. It would also be of particular interest to develop novel compounds that exhibit a favorable therapeutic profile in vivo (e.g., are safe and effective, while retaining stability in biological media)."}
{"text": "The science of transplantation, now half of a century old, has dramatically increased and improved the life of many individuals, including many children, with end stage diseases. Recent advancements in immunosuppressive agents have substantially decreased rejection of allografts over the past decade and a half in the United States. Hariharan et al., N Engl J Med 342, 605-612 (2000); and First, M. R., Transplant Proc 34:1369-1371 (2002). However, to avoid both episodes of acute rejection and the initiation of chronic rejection following transplantation, immunosuppressive drugs must be administered over the entire life of the organ recipient. Consequences of this long-term administration are profound, including undesirable side effects, increasing the risk of infection, autoimmunity, heart disease, diabetes, and cancer. The chronic administration of these immunosuppressive drugs (especially when give systemically) lead to toxicity and significant side effects, thereby leaving the patient vulnerable to a variety of diseases and systemic organ failure.\nThe most desirable alternative to this extended state of vulnerability would be to render the patient's immune system to effectively suppress immune activation without systemic immunosuppression. In this case, no further immunosuppressant drug treatment would be necessary. Furthermore, the recipient's immune system would otherwise function normally, being capable of combating pathogens and malignant tumor cells. What is needed in the art are compositions and methods to induce permanent immunological tolerance by stimulating immunosuppressive regulatory cells."}
{"text": "(a) Field of the Invention\nThis invention relates to an alarm system for a machine tool and more particularly but not by way of limitation to a portable alarm system used with a CNC milling machine for alerting a machine operator that the milling of a particular work piece has been completed and the milling machine needs to be attended to and a new work cycle started.\n(b) Discussion of Prior Art\nHeretofore, there have been a number of patent applications filed on various types of audio and visual alarms used with milling machines for sounding an alarm related to a cutting fluid shortage, overheating of a machine tool or work piece, collisions of work pieces and moving of axes and tools. Also, in U.S. Pat. No. 6,717,514 to Stein et al., U.S. Pat. No. 5,579,884 to Appleyard et al, U.S. Pat. No. 5,563,581 to Kats, U.S. Pat. No. 5,486,691 to Dieterle, U.S. Pat. No. 5,148,053 to Dubois III, and U.S. Pat. No. 4,973,944 to Maletta, various types of machine tool safety devices, monitoring apparatus, and alarm devices are described for preventing operator injuries when using lathes, milling equipment, metal shears and the like.\nNone of the above mentioned prior art alarms have the particular unique features, objects and advantages of the subject portable alarm system used with a tool magazine and a CNC milling machine spindle."}
{"text": "The present invention relates to a white light-emitting diode device which has an excellent color rendering property and which can be used in a back light and a front light as an externally irradiating light source of a color liquid crystal display and for room illumination. Particularly, the present invention relates to a covering member that enables a conventional blue or green light-emitting diode device to emit white light, and a white light-emitting diode device using the covering member.\nA light-emitting diode (hereinafter, referred to as xe2x80x9cLEDxe2x80x9d) emits a light when electrons from an N-region and holes from a P-region are moved to a PN junction section and are recoupled in a PN junction semiconductor when a forward voltage is applied. Since the energy released when the free electrons are coupled is radiated as light, the light emitted from the LED is light in a narrow wavelength range, i.e. light of a single color such as red or blue.\nOn the other hand, a light source for illumination of a color liquid crystal display and for room illumination needs to emit fluorescent light or white light approximate to natural light.\nThus, a full-color LED device has been used as a light source for illumination of the color liquid crystal display and for room illumination. The full-color LED device generates white light by using a combination of a red light-emitting element for emitting red light, a blue light-emitting element for emitting blue light and a green light-emitting element for emitting green light.\nHowever, with a full-color LED device as a light source results in a complicated design, since each of the light-emitting elements in the full-color LED device requires must have an anode terminal and a cathode terminal. Further, a larter space and higher costs are required since a full-color LED device uses three kinds of light-emitting elements. In addition, if one light-emitting element of a full-color LED device is broken, the desired color balance is lost, with the result that the intended color tone may not be obtained. In view of the above situation, many studies are being vigorously undertaken to develop an LED device capable of emitting white color using one kind of light-emitting element. Such an LED device is called a white LED device and is distinguishable from the full-color LED device. Using a white LED device enables a back light or the like to be driven by a simple electric circuit without requiring an inverter circuit and the like, and reduces power consumption. This presents advantages such as reducing the outer configuration of the driving circuit and eliminating electromagnetic noise. Because of such advantages, the white LED device is being developed.\nSince a white LED device must generate white light by using only one kind of light-emitting element (single-color), it has been the general practice to use a single-color light-emitting element in combination with a phosphor that can convert the wavelength of light emitted from the light-emitting element to emit light of another color. For example, white LED devices which have a layer containing YAG (yttrium aluminate) phosphor provided on a blue light-emitting element are commercially available.\nThe YAG phosphor generally emits light in the green to yellow region, depending on the kind of activating element in the YAG. Thus, the red light obtained in a white LED by combining a blue light-emitting element for emitting light having an excitation wavelength of YAG and the YAG phosphor is insufficient. This leads to the problem that matters to be displayed in red look subdued.\nTo remedy this problem, it has been considered to use a YAG phosphor capable of emitting fluorescent light approximate red light, such as by activating YAG, for example, by substituting Gd at the Y-site of Y3Al5O12, to activate the YAG phosphor and shift the emission peak toward a longer wavelength, i.e. toward the red side of the spectrum. However, such a method cannot be used to provide a white LED device having an average color rendering evaluation number (Ra) of 85 or larger, the value of Ra being an index as to whether an illuminated object conforms to color perception under a reference light. Although more red light may be obtained by increasing the content of red YAG in the fluorescent layer, the emission intensity of yellowish light emitted from YAG also becomes higher, with the result that an average color rendering evaluation number Ra of 85 or larger cannot be obtained.\nAs a method of compensating for not only the blue, green and yellow lights, but also the red light in the white LED, the use of a red phosphor has also been considered. However, at present the manganese activated magnesium arsenate phosphor known as a dark red phosphor is not being produced for fluorescent lamps due to its toxicity. Further, since manganese activated magnesium fluoride germanium oxide phosphor is expensive, it is not suitable for the use of a white LED as a back light having the merits of being light-weight and thin and of requiring low power consumption.\nThe effect of adding red pigment to the fluorescent layer of an LED is also being studied (see, for example, Japanese Unexamined Patent Publication No. (HEI)5-152609). However, since the pigment acts like a light filter in making a specified color stand out by cutting other colors, addition of the pigment results in a reduced luminance of the LED device and, thus, is not suitable for a white LED device aimed at conserving space, costs and power.\nAn object of the present invention is to provide, without increasing manufacturing costs more than necessary, an LED device having an excellent color rendering property and a covering member used for such an LED device.\nAn inventive light-emitting diode comprises a light-emitting element for emitting light in the blue to green region and a fluorescent substance containing a CaS phosphor activated by Eu for converting the wavelength of the light emitted from the light-emitting element to a different wavelength.\nAnother inventive light-emitting diode device comprises a light-emitting element for emitting light in the blue to green region and a fluorescent substance containing a phosphor expressed by a general formula AEu(1xe2x88x92x)LnxB2O8 for converting the wavelength of the light emitted from the light-emitting element to a different wavelength, wherein A is an element selected from the group consisting of Li, K, Na and Ag; Ln is an element selected from the group consisting of Y, La and Gd; and B is W or Mo; and x is number equal to or larger than 0, but smaller than 1.\nAn inventive covering member for a light-emitting diode device which includes a light-emitting element sealed by a transparent resin for emitting light in the blue to green region comprises a base polymer and a CaS phosphor activated by Eu.\nAnother inventive covering member for covering a light-emitting diode device which includes a light-emitting element sealed by a transparent resin for emitting light in the blue to green region comprises a base polymer, and a phosphor expressed by a general formula AEu(1xe2x88x92x)LnxB2O8, wherein A is an element selected from the group consisting of Li, K, Na and Ag; Ln is an element selected from the group consisting of Y, La and Gd; and B is W or Mo; and x is number equal to or larger than 0, but smaller than 1.\nThe inventive light-emitting diode devices are better in emitting red light than conventional light-emitting diode devices using a light-emitting element that emits light in the blue to green region. Thus, according to the present invention, a white diode device having a color rendering property approximate to that of a fluorescent lamp or of natural light can be provided by containing YAG phosphor in the fluorescent substance.\nFurther, if the inventive covering member is mounted on a conventional light-emitting diode device, it can be used as the inventive white LED device having an excellent color rendering property.\nThese and other objects, features and advantages of the present invention will become more apparent upon a reading of the following detailed description and accompanying drawings."}
{"text": "1. Field of the Invention\nThis specification relates to a refrigerator having a vegetable storage container, and particularly, to a hermetic structure of a vegetable storage container for a refrigerator capable of sealing an inner accommodation space of a vegetable box.\n2. Background of the Invention\nIn general, a refrigerator is equipment for keeping food stuffs in a cold state for a long term of time by generating cold air by way of driving a refrigerating cycle installed in the refrigerator and supplying the cold air into a refrigerating chamber and a freezing chamber.\nThe refrigerator may include a vegetable container for keeping vegetables. When the vegetables are kept in the refrigerator, they should be maintained as fresh as possible. Therefore, it may be important to maintain the space for storing the vegetables in an optimal condition.\nFIGS. 1 to 3 show main parts of a vegetable container for a refrigerator according to the related art, which are a perspective view and coupled views of a sealing structure of the vegetable container for the refrigerator according to the related art.\nReferring to FIG. 1, the related art vegetable container may include a vegetable box 200 having an accommodation space having an upper opening for storing vegetables therein, and a cover 100 for covering the upper opening of the vegetable box 200.\nThe cover 100 may include an inner cover 120 having cover gaskets 123 for sealing the upper opening of the vegetable box 200, and an outer cover 110 located on the inner cover 120.\nThe vegetable box 200 may include a main body 210 defining the accommodation space for the vegetables therein, and a front part 230 having a front panel 231 formed thereon.\nThe cover 200, as shown in FIG. 2, may be pushed in a slidable manner to close the upper opening of the accommodation space of the vegetable box 200.\nEspecially, cover rails 113 may be disposed long in back and forth directions on both sides of a lower surface of the outer cover 110 of the cover 100. The cover rails 113 of the outer cover 110 and main body rails 211, which are disposed on both sides of an upper surface of the main body 210, may be mutually slidable to open or close the inner accommodation space of the vegetable box 200.\nReferring to FIG. 3, the cover 100 may be closed in a manner that stopping protrusions 1131 of the outer cover 110 are engaged with protrusions 2111 formed on end portions of the rails 211 on the upper surface of the main body 210 such that the cover 100 cannot be slid.\nAccordingly, when the cover 100 is closed with sealing the upper opening of the vegetable box 200, the cover gaskets 123 may be pressed due to the weight of the cover 100, thereby realizing a sealing function.\nHowever, the vegetable box for the refrigerator according to the related art may have the following problems.\nReferring to FIG. 3, the related art vegetable box cover 100 may cover the upper opening of the vegetable box 200 in the sliding manner. However, without a force applied upwardly, the cover 100 and the vegetable box 200 may be separated from each other by a predetermined gap at a gasket-loosed portion S. This may make it difficult to completely engage the cover gaskets with the vegetable box, and accordingly a gap may be generated. Also, the accommodation space of the vegetable box 200 may not be completely sealed by the cover 100.\nConsequently, odors of food stuffs stored in the refrigerating chamber or the freezing chamber of the refrigerator may permeate vegetables stored in the inner accommodation space of the vegetable box 200 via the gasket-loosed portion S. Also, the accommodation space of the vegetable box 200 may not be maintained with proper humidity, and accordingly the vegetables may not be kept in a fresh state.\nAlso, due to the loosed gap, the vegetable box may not be sealed depending on a user's usage pattern. Accordingly, a sealing force shortage may be caused and the vegetables may be caught."}
{"text": "In ventilation systems, the ventilation air is admitted to the room through openings in the wall or ceiling and such openings are fitted with a removable ventilation elements, normally air valves, diffusers or grilles for supplying and exhausting air. The conduits to which the elements are connected may be of varying widths or diameters and may even vary in type. Frequently the conduit is simply an opening in a concrete wall and such conduits are often uneven and include a great number of varying irregularities.\nIf the conduit embodies a standard piece of duct work or other pre-formed conduit, the ventilation element may be designed to mate with the pre-formed conduit and may be screwed in place or glued by suitable sealing cement. However, where it is desirable to permit ready removal and replacement of the ventilation element, it has been proposed to provide spring latches which releasably mount the ventilation element in the duct. Such devices are normally not suitable unless the ventilation element has a cylindrical portion which telescopically engages in the duct conduit to position the same. In such cases, the latch element may comprise a special expanding member on the cylindrical portion which springs out to engage the conduit and secure it in position.\nIn order to solve these problems, a mounting arrangement has been provided in which the ventilation element is provided with resilient brackets which may spring outwardly against or behind an irregularity or a shoulder. The resilient element then is designed either to pass \"over-center\", or to act as a cam element to resiliently cause the ventilation element to bear tightly against the wall surface surrounding the ventilation opening.\nThe foregoing mounting apparatuses are not entirely satisfactory, because of their complicated construction and assembly, their expense, and their inability to adapt to openings which are quite irregular."}
{"text": "Generally speaking, conventional lidding films either: 1. provide reliable seal integrity but are often difficult to open; or 2. are easily opened but provide less reliable seal integrity. The present invention eliminates this trade-off; that is, the lidding films of the present invention are easily opened by consumers and have reliable heat seals, even at very low temperatures.\nU.S. Pat. No. 4,375,494 to Stokes discloses multilayer films having a semi-crystalline polyester support layer and an amorphous polyester heat seal layer. The heat seal layer further comprises a finely divided particulate additive."}
{"text": "Patent Document 1 discloses a semiconductor manufacturing device. The semiconductor manufacturing device described in Patent Document 1 includes a stage and a plurality of pins. The stage has a plurality of holes into which the pins are respectively inserted.\nSome of the pins are used for vertically moving a substrate to be processed, which is mounted on the stage, with respect to the top surface of the stage. The other the pins are used for vertically moving a focus ring mounted on the stage with respect to the top surface of the stage.\nPatent Document 1: Japanese Patent Application Publication No. 2006-196691\nIn the above semiconductor manufacturing device, a plurality of driving units for vertically moving the pins is required as many as the number of the pins, so that the device becomes complicated. Further, in the above semiconductor manufacturing device, holes as many as the number of pins need to be formed in the stage. However, it is preferable to form a smaller number of holes in the stage."}
{"text": "Fiber to the x (FTTx) is a generic term for broadband network architecture that uses optical fiber to replace the copper cabling usually used for connecting the network to the end user. E.g. fiber-to-the-building (FTTB) is a term that defines the scenario when fiber reaches the boundary of the building, such as the curb or basement in a commercial or residential building. The fiber is there connected to some kind of connection device where the final connection to the tenants is made via alternative means.\nAn Optical Network Unit (ONU) is an example of such a connection device. When installed close to, or in, a commercial or residential building the connection device is often referred to as a Multi-Dwelling Unit (MDU). The ONU delivers services over Very-high-speed Digital Subscriber Line (VDSL2) or Ethernet from a Giga-bit or Ethernet Passive Optical Network (GPON or EPON) feed. The services are provided to the ONU via a fiber connection coming from an access node (OLT—Optical Line Termination) in a central office. The ONU converts the optical signals in the optical fiber to electrical signals for delivery to the tenants. The ONU is not powered locally, but retrieves power from a remote power station in a remote terminal (RT) that can be located far away. The powering is provided via already existing twisted pair telephone copper from the RT to the ONU. The twisted pair telephone copper was once being used for Plain Old Telephone Service (POTS) and later for Asymmetric Digital Subscriber Line (ADSL) containing POTS. Due to the small wire gauge of the twisted pair telephone copper the power supply in the RT must deliver the power using high DC voltage since the use of low DC or AC voltage would result in high losses. The RT may therefore up-convert −48V (which is the typical operating voltage of the RT) to +/−190V which is necessary to deliver enough power (˜95 W per power channel, 380V, 250 mA max) to the ONU. Depending on power consumption and distance, one or two power pairs are typically utilized. Accordingly, the ONU down-converts +/−190V to −48V to power internal electronics. The RT may be placed in the central office or in another location where access to the AC power grid is available.\nTo protect the RT and ONU equipment from surges during lightning, primary (gas discharge tubes) and secondary (sidactors) protection mechanisms are deployed in both the RT and the ONU, effectively shortening power leads to ground, bleeding surge energy off. This causes power outages in the area of tens of milliseconds that leads to service unavailability for users. Therefore, hold-over capability in form of an inexpensive hold-over circuit containing low-lost and low-value capacitors may be installed at the ONU to power the device during the short time period of power loss from the remote up-converter.\nWhen a ground fault is detected on the power leads, the up-converter performs a safety power shutdown of the RT for between 60 milliseconds to several seconds according to safety standards. This results in loss of power from the up-converter of the RT to the ONU. For economic reasons an ONU should preferably not be provided with any expensive hold-over capability of several seconds such as ultra-capacitors, batteries, etc. Consequently, a ground fault leads to a reboot of the ONU leading to service outage to all costumers of the ONU for some minutes.\nGround faults do not occur frequently since they are introduced by safety breaches caused by humans or animals touching the leads or power crossover. However, the poorly shielded and aged telephone copper that feeds from RT to the ONU can be very sensitive to external disturbances such as interferences from close-by power distribution and transmission lines causing common mode transients which result in current imbalance. Sometimes the up-converter misinterprets these disturbances as ground fault events causing unnecessary safety power shut downs.\nThere is thus a need for a method and an apparatus that is resilient to interferences but still complies with existing safety standards."}
{"text": "Absorbent articles for personal hygiene of the type indicated above are designed to absorb and contain body exudates, in particular large quantity of urine. These absorbent articles comprise several layers providing different functions, for example a topsheet, a backsheet and in-between an absorbent core, among other layers. The function of the absorbent core is typically to absorb and retain the exudates for a prolonged amount of time, minimize re-wet to keep the wearer dry and avoid soiling of clothes or bed sheets.\nThe majority of currently marketed absorbent articles comprise as absorbent material a blend of comminuted wood pulp with superabsorbent polymers (SAP) in particulate form, also called absorbent gelling materials (AGM), see for example U.S. Pat. No. 5,151,092 (Buell). Absorbent articles having a core consisting essentially of SAP as absorbent material (so called “airfelt-free” cores) have also been proposed (see e.g. WO2008/155699 (Hundorf), WO95/11652 (Tanzer), WO2012/052172 (Van Malderen)). Absorbent cores with slits or grooves have also been proposed, typically to increase the fluid acquisition properties of the core or to act as a folding guide.\nWO2012/170778 (Rosati et al., see also WO2012/170779, WO2012/170781 and WO2012/170808) discloses absorbent structures that comprise superabsorbent polymers, optionally a cellulosic material, and at least a pair of substantially longitudinally extending channels. The core wrap can be adhesively bonded through the channels to form a channel bond. The channel bonds may be permanent, so that their integrity is at least partially maintained both in dry and wet state. As the absorbent structure absorbs liquid and swells, the absorbent structure takes a three-dimensional shape with the channels becoming visible. The channels are indicated to provide improved fit and/or better liquid acquisition/transportation, and/or improved performance throughout the use of the absorbent structure. Any superabsorbent polymer particles known from the superabsorbent literature are indicated to be suitable.\nThe properties of superabsorbent polymers have been characterized in various ways. The absorbent capacity (CRC) in grams of liquid per gram of superabsorbent particles has been used, as well as their absorption speed as measured by the Free Swell Rate (FSR) and their permeability as measured by the Urine Permeability Measurement (UPM) test.\nInternational patent application WO2012/174,026A1 discloses the K(t) method which can be used to determine the time dependent effective permeability (K(t)) and the uptake kinetics (T20) of a gel layer formed from hydrogel-forming superabsorbent polymer particles under a confining pressure. The application indicates that these SAP can be used to reduce leakage, especially at the first gush, i.e. when the article starts to be wetted.\nIt has now been found that although the absorption properties of conventional SAP may not be negatively impacted at first gush when used in a core with channels, the liquid absorption of the SAP can be significantly reduced in the following gushes after the fluid has been already absorbed in these cores comprising channels compared to cores without channels. Without wishing to be bound by theory, the inventors believe that the three-dimensional channels which are formed as the SAP absorbs a fluid can create a resistance to swelling for the superabsorbent polymers and reduce their swelling kinetics. As the channels otherwise facilitate the distribution of the fluid along the core, it was on contrary expected that any conventional SAP could be used in these cores. Accordingly the inventors have found that for absorbent cores comprising such channels it can be advantageous to use these SAP having a T20 of below 240 s to maintain sufficient speed of absorption beyond first gush."}
{"text": "The present invention relates to the identification of articles, especially golf balls, which are marked, distributed from a base and subsequently collected.\nGolf driving ranges are becoming increasingly popular, especially those involving scoring and other competitions. In such cases the ability to identify a collected ball is important Even in conventional driving ranges, the ability to identify golf bell can facilitate charging players and or detecting theft.\nU.S. Pat. No. 5,439,224 discloses a golf range with targets provided with optical scanners to identify balls and to pass the information to a programmed computer. U.S. Pat. No. 5,370,389 discloses a golfing range game in which a ball rolling into a hole is sensed by a scanner which identifies the tee from which the ball was driven by detecting the colour of the ball or a bar code printed thereon; the information is passed back to a display near the golfer. U.S. Pat. No. 5,653,642 discloses a driving range with targets with optical readers for detecting unique identification ball codes, the optical readers being connected to a computer device.\nSystems employing optical codes have a number of disadvantages. In particular, dirt on the ball can prevent correct identification. Also, the optical codes are subject to wear and abrasion. Furthermore, the orientation of the ball relative to the optical reader is important in order for the code to be correctly read. In addition the coding markings can be visually intrusive.\nGB-A-2267222 discloses a golf driving range employing balls which are distinguished from each other by means of bar codes or transponders. Ball detecting means are provided at the holes and the surrounding greens. The disclosure of this document corresponds to the introductory part of claim 1.\nU.S. Pat. No. 5,626,531 discloses golf balls with passive tags which can be sensed by an electronic detection system and U.S. Pat. No. 5,743,815 discloses a golf ball housing a passive transponder which may be uniquely coded.\nU.S. Pat. No. 4,516,770 discloses a ball detection and identification arrangement for a table ball game, e.g. pool. Balls landing in the pockets pass to an electronic detector which operates a display and scoring device.\nU.S. Pat. No. 5,487,542 discloses a golf game with identifiable golf clubs and means for sensing the presence of a golf ball in a hole. In a modification the balls may be uniquely coded optically.\nThe present invention seeks to overcome or reduce the above disadvantages.\nAccording to a first aspect of the present invention there is provided a golf driving range employing uniquely-coded golf balls characterised in that a tee of the range incorporates ball-identifying means connected to a computer system so that the computer system knows which ball a player is driving from said tee.\nAccording to a second aspect of the present invention there is provided a golf driving range comprising means for supplying golf balls to a player, said supply means incorporating first ball-identifying means, the range further comprising means for collecting driven balls incorporating second ball-identifying means, the range further comprising means for processing the output of said first and second ball-identifying means, wherein each golf ball incorporates a coded RF-transponder and the first and second ball-identifying means employ RFID technology.\nAccording to a third aspect of the present invention, there is provided a golf driving range comprising means for supplying golf balls one-by-one to a player, said supply means incorporating first ball-identifying means, the range further comprising means for collecting balls driven by the player, said collecting means incorporating second-ball identifying means, the range further comprising means for processing the output of said first and second ball-identifying means, each golf ball incorporating a RF-transponder having a unique code, and the first and second ball-identifying means employing RFID technology, wherein the golf balls to be supplied to the player are fed one-by-one from a holder of a plurality of golf balls past the first ball-identifying means and are then supplied in the same order to the player.\nThe outputs of said first and second ball-identifying means are preferably supplied to a computer, which processes the received information and provides a display of the result.\nThe supplying and collecting means preferably comprise channels for directing golf balls, and the first and second ball-identifying means are preferably RFID readers with reading heads located directly adjacent the respective channels.\nThe balls collected by the collecting means may be arranged to be supplied to the supplying means for recycling.\nAccording to a fourth aspect of the present invention there is provided an apparatus for supplying golf balls one-by-one to a player and incorporating ball-identifying means and means for processing the output of said ball-identifying means, each golf ball incorporating a RF-transponder having a unique code and the ball-identifying means employing RFID technology wherein the golf balls to be supplied to the player are fed one-by-one from a holder of a plurality of golf balls past the ball-identifying means and are then supplied in the same order to the player."}
{"text": "The present invention relates in general to a method of forming a protective coating on a substrate by means of plasma polymerization with at least one gaseous silicon-organic compound, in which a plasma is generated in a vacuum chamber and a gas with the compound is fed thereto for the production of the coating. The invention relates further to a coating apparatus with a coating chamber, a source arrangement which, for the generation of a plasma, generates an electromagnetic alternating field, a pumping arrangement connected to the chamber, and a controlled gas inlet arrangement for the infeed of gas into the chamber.\nIn the technology of protective coatings it is generally known to provide substrates, specifically non-planar substrates such as reflectors for automobile or other headlights, floodlights, etc. with corrosion proof or wiping resistant, protective coatings.\nPlasma polymerization which, as is generally known, is used for this task, and which is a plasma enhanced chemical deposition method operative at comparatively low substrate temperatures, permits metal substrates or previously coated metal layers, specifically aluminum vapor coated substrates, to be coated with siloxane, a silicon-organic compound.\nBecause a plasma is generated for this procedure by direct current, as disclosed, for example, in the German patent publication DE-OS 22 63 480 or French patent reference FR A 21 690 072, a new metal encasement must be furnished from time to time, around at least one of the metal electrodes used for the coating process and between which the plasma is generated after every coating process. This is because an insulating coating is produced at relatively high coating rates on the electrode during coating with siloxane. Without this new metal encasement, the longer the process continues, the more dielectric coating is formed on the electrode. This would lead to an unstable plasma with flashovers through the insulating coating and contamination of the coating process area, sputtering-off until, finally, a disruption or break down of the plasma discharge takes place.\nThis problem remains if, according to German reference DD A 272 773, a pulsed DC-plasma is used. According to DD A 272 773, an alternating voltage of selected frequency is rectified in one direction and is used for the plasma generator. By means of a super-imposed DC-voltage, the total voltage is raised causing the plasma discharge to intensify which, it is said, increases productivity. It is also stated in the reference that an upper limit for this DC-voltage exists, above which the discharge becomes continuous. This pulsed switching of a discharge with a non-pulsed operating voltage (a rectified voltage) can no longer support pulsed DC-discharge operation. A pulsed DC-discharge is used according to DD A 272 773 for preventing arcing, which, as is known, is often encountered in pure DC-discharge operation.\nDue to the necessity of having to reencapsulate the mentioned electrode, relatively long maintenance intervals occur due to which a production plant must often shut down for prolonged time spans. This can not be tolerated with in-line production plants, for example.\nIt is noted that the parameters for a plasma polymerization method depend very strongly on the gaseous compounds used and which will thus form the process coating. It is known from European patent reference EP A 207 767, to pulsate RF plasma, during a surface treatment-etching or coating by plasma enhanced reactive processes. Without any specific selection, a large number of different materials are proposed for processing in this reference, e.g., Si.sub.3 N.sub.4, TiO.sub.2, Al.sub.2 O.sub.3, BN, SiO.sub.2, B.sub.4 C, SiC, HC, TiC, TiN, BP. All of these coating materials are not produced by polymerization.\nIn case of the plasma polymerization of silicon-organic compounds, it is known, e.g. from EP A 0 299 754 to generate the plasma by means of an electromagnetic alternating field. It is further known that, when using certain non-silicon organic gaseous compounds such as methane, higher deposition rates can be obtained by pulsing the electromagnetic alternating field, and at the same applying average power.\nHow the coating process behaves, as noted above, depends on plasma modulation, and to a large extend, on the specific gaseous compound supplied for the coating operation.\nRegarding the recognition that it is possible to reach, in certain cases, higher coating rates by a pulsed electromagnetic alternating field for the generation of the plasma at the same average electrical power applied also under the same process conditions, reference is made to:\nVinzant J., Shen M., Bell A. ACS SYMP.SER; SHEN BELL 108 p. 79 (1979), \"Polymerization of Hydrocarbons in a pulsed plasma\"; PA1 H. Yasuda, T. Hsu: J. Polym Sci, Polym Chem ed 15, 81, (1977), or PA1 H. Yasuda \"Plasma polymerization\" Academic press (1985), p.103-105; and PA1 Lloret A., Bertran E., Andujar J., Canillas A., Morenza J., J.APPL.PHYS. 69, p. 632 (1991), \"Ellipsometric study of a-Si:H thin films deposited by square wave modulated rf glow discharge\". PA1 GB A 2 105 729 (1981), ITT Industries Ltd., London-GB; and PA1 U.S. Pat. No. 4,950,956 (1990), Anelva Corp., Tokyo, Japan. PA1 JP 62 103 371 (1985), Hitachi KK; PA1 DD 264 344 A3 (1986), VEB ZFT Mikroelectronik. PA1 in relation to the depositing of diamond or graphite out of methane, from JP 62 123 096 (1985), Showa Denko KK; PA1 for silane (SiH.sub.4) from: Watanabe Y., Shiratani M., Kubo Y, Ogawa I., Ogi S., APPL.PHYS.LETT. 53, --p. 1263 (1988), \"Effects of low frequency modulation on rf-discharge chemical vapor deposition\"; PA1 with reference to the deposition of amorphous carbon from methane, from Hossary F., Fabian D., Webb A., THIN SOL.FILM 192, p.201 (1990), \"Optical properties of amorphous carbon films formed by rf-plasma deposition from methane\"; and PA1 again with reference to silane (SiH.sub.4) further reference is made to: Watanabe Y., Shiratani M., Makino H., APPL.PHYS.LETT. 57, p. 1616 (1990), \"Powder-free plasma chemical vapor deposition of hydrogenated amorphous silicon with high rf power density using modulated rf discharge.\"\nFurther, that a lower heat loading of the substrate takes place by the application of a pulsed electromagnetic alternating field for the generation of the plasma follows from:\nIt is furthermore, also known that new coating properties are realized by the application of a pulsed electromagnetic alternating field for the generation of a plasma. Reference is made in this respect to:\nImprovement in the purity of the layer or a reduction of powdery layer portions is known:\nAll of the above documents are to be taken as an integrated part of the present description and are to be considered as incorporated in the present description by reference."}
{"text": "    [Patent document 1] JP 2008-205604 A\nAn antenna is mounted in a vehicle for executing various communications such as GPS, vehicle-to-vehicle communications, or road-to-vehicle communications. Such an antenna in a vehicle is required to enable downsizing and favorable transmission and reception performance.\nSeveral wireless communication modules with antennas are provided for the several kinds of communications in a vehicle, respectively; the several wireless modules are required to be arranged collectively at the same place (for example, a predetermined place on the roof of the vehicle). Therefore, the various wireless communication modules with the antennas are required to achieve further downsizing.\nTo that end, Patent document 1 discloses a technology to downsize a wireless communication module with an antenna as follows. A radio frequency circuit board is mounted on a ground structure serving as a ground of a wireless terminal. A circuit board has a feed point at one end and a variable capacity reactive element at the other end. The feed point has one end connected with one end of a ground conductor of the circuit board and the other end connected with the ground structure via a first contact terminal. The variable capacity reactive element has one end connected with the other end of the ground conductor of the circuit board and the other end connected with the ground structure via a second contact terminal.\nThis forms a current loop that flows from the one end of the feed point, via the first the ground structure, the second contact terminal, the variable capacity reactive element, and the earth conductor, to the other end of the feed point. This configuration achieves a radiating pattern shape as an antenna without using an antenna element, thereby enabling downsizing.\nHowever, the feed point and the variable capacity reactive element are arranged at both the ends of the circuit board, respectively, to achieve a function as one antenna; thus, the downsizing is not fully achieved.\nFurther, in the antenna apparatus of Patent document 1, the above-mentioned current loop is formed, thereby causing both the ground structure and the earth conductor to be radiation source. This requires the design to consider the both. Therefore, the design of the antenna apparatus may become complicated; the configuration of the antenna apparatus may become complicated."}
{"text": "This invention relates to the construction of an electronic watch, and particularly to the manner in which a watch substrate is mounted, protected and controlled within a watch case.\nAn electronic watch is one in which time increments are generated at a frequency in the order of kilocycles to megacycles per second, with dividers and memories for electronic time information processing. Readout is conveniently digital so that there are usually no movable mechanical parts except for switches to control the electronics. In recent years, there have been many inventions in electronic watches and many of the patents describe at least a portion of the physical constructions which interrelates the physical and electronic components. Pertinent background patents include McCullough et al. U.S. Pat. No. 3,759,031, Perkins et al. U.S. Pat. No. 3,784,725, Doss U.S. Pat. No. 3,846,972, Yamazaki U.S. Pat. No. 3,800,523 and Zurcher et al. U.S. Pat. No. 3,838,567. These patents are just a few in the large body of prior art in the electronic watch and related field and are offered as examples of patents which disclose some of the physical structure of electronic watch construction.\nIn assemblying an electronic watch the electronics must be supported and protected and must be related to other components for physical and electrical interconnections. The clamping, attachment, interconnection, and protection of the watch structure is a field in which concepts relating to economy reliability and serviceability must be applied. The prior art does not indicate the manner in which various design features should be optimized."}
{"text": "The present invention relates to a wireless thermostatic control system based on room occupancy. More particularly, the present invention relates to a wireless thermostat control system based on room occupancy for saving energy by detecting whether or not a person is in the room and setting back energy requirements if the room is not occupied. In the hotel and motel industry, the cost of heating and cooling a room has been of concern for a long time. For example, U.S. Pat. No. 4,315,596 issued in 1980 to Johnson Jr. et al. discloses an energy conservation system for inns, hotels and motels for reducing the unproductive consumption of energy, primarily during periods where a room is unoccupied. Energy consumption during periods of non-occupancy is also of concern to individual consumers and in office buildings and other structures. Other patents of interest in this field include U.S. Pat. Nos. 5,127,575 to Beerbaum, 5,165,465 to Kenet, and 4,101,886 issued in 1978 to Grimes et al. All of these systems are complicated and require substantial hard wiring to install. There are also commercially available systems such as Linear Corporation's 9HC40 and GEM 5000 and North American Technologies'Sensorstat.RTM.1000, 2000 and 3000 which suffer from the same problems of complexity and difficulty of installation as the above patents.\nFor example, in the Grimes, et al. patent, the system includes a door switch, a dead bolt switch, a transmitter, a receiving means for receiving signals transmitted from the transmitter, a connecting means for connecting the receiving means to the particular appliance, such as a heating and cooling system's thermostat, and there is hard wiring not only between the door switch, dead bolt switch and transmitter and but also between the receiving means and appliance. The system will, upon the dead bolt being opened and the door being opened and closed, cause the transmitter to send out a signal to the receiving means connected to the existing controls of the particular appliance. There is a two and a half minute delay until the appliance enters a set back or turned off mode. When the room is reentered, the door switch detects opening and closing of the door. Reengagement of the dead bolt will then cause the transmitter to send out a signal to fully activate the appliances. Where set back control is used, the receiving means contains the extra thermostats and control circuitry to perform this function. The receiving means must thus have a separate source of energy from the appliance being controlled and must be hardwired to the existing thermostat. The commercially available systems use similar operation concepts as Grimes et al., but also use sensors to detect whether or not someone is in the room. In fact, these systems have hard wiring between a door switch and a power pack and between a sensor and the power pack associated with the extra thermostat control unit."}
{"text": "1. Field of the Art\nThis invention relates to niobium-containing weldable structural steel having good weldability. More specifically, the present invention relates to a niobium-containing weldable structural steel having yield strength of 40-70 kg/mm.sup.2 prepared by providing a niobium-containing steel with a predetermined composition so as to improve toughness in the weld heat affected zone and resistance to weld cracking.\n2. Description of the Prior Art\nWhen added in a small amount to steel, niobium improves the strength and toughness of steel, and is an economically advantageous element. It is for this reason that a niobium-containing steel (hereinafter referred to as \"Nb-containing steel\") has gained a wide application as a weldable structural steel for pipe lines, ship-building, pressure containers, bridges and the like. Among numerous applications of the Nb-containing steel, the detailed description is hereby presented about line pipe steels used in specifically great quantities. Nb-containing non-quenched tempered high tensile strength steel has conventionally been used in great quantities for pipe lines for transporting crude oil and natural gas. However, the pipe lines that have so far been laid down have a small pipe diameter and low inner pressure and are used at a relatively high temperature of not lower than about 0.degree. C. Hence, requirement for toughness in the weld heat affected zone has not been much serious.\nAs laying of pipe lines has gradually been concentrated on the ultra-cold region such as in U.S.S.R., Canada, U.S.A., etc. and the material transported has also been changed gradually from crude oil to natural gas, however, an unexpected problem has come to the front. The conventional Nb-containing steel for pipe line causes embrittlement of the weld heat affected zone (HAZ) in the seam weld portion at the time of steel making and hence, fails to ensure sufficient toughness to withstand the use in such regions. As countermeasures, a multiple-electrode submerged arc welding process of a low restricted weld heat input or a MIG welding process has been examined on one hand in the aspect of the welding process, in order to mitigate the influence of the welding heat. Research and development has been made on the other hand in the aspect of the steel material to be used, in order to obtain a steel having such a composition of economical components that does not cause degradation of the weld heat affected zone. With regard to the properties of the steel material especially, the most desirable steel would be one that does not cause embrittlement of the weld heat affected zone and yet exhibits good notch toughness even when welding of a large weld heat input such as one side welding or both side welding with one or two layers (not more than three passes in each weld groove) is carried out in order to improve the welding efficiency and reduce the cost of production of the pipe. At the same time, the steel must satisfy a specific requirement of low weld crack sensitivity for the purpose of preventing cracks of the weld portion because when the on-the-site welding is effected in the cold zone, a high cellulose type electrode of a high hydrogen content is generally employed without preheating of the pipe in the low temperature atmosphere.\nIt has been a pressing need for those concerned in the art to develop a steel strip for a line pipe which simultaneously satisfies the abovementioned requirements, i.e., to have excellent toughness in the weld heat affected zone and to be less hardenable and excellent in its resistance to weld cracking.\nIncidentally, the following U.S. patents are located as the prior art most relevant to the present invention;\nU.S. Pat. Nos. 3,592,633, 3,807,990, 3,619,303, 3,725,049, 3,721,587 and 3,132,025."}
{"text": "This invention relates to a propulsion unit for an inboard-outboard motor and more particularly to an improved propulsion unit that lends itself to reversal of the rotation of the propulsion device in a simple and expedient manner for twin outboard drive arrangements.\nIt is well known that a marine outboard drive generally has a side thrust due to the direction of rotation of either the propeller or other propulsion device. This side thrust may be countered by employing twin counterrotating outboard drives. For this and a variety of other reasons, it has been increasing practice to use such twin outboard drives, be they outboard motors or the outboard drive unit of an inboard-outboard arrangement.\nIn connection with such twin outboard drives, normally there is a separate internal combustion engine that powers each outboard driven. In order to achieve the counterrotation of the outboard drives, either the internal combustion engines must rotate in opposite directions or the gearing associated in the drive between the engine and the propulsion device must include a reversing mechanism for reversing the direction of rotation. In addition, each outboard drive normally includes its own forward, neutral, reverse transmission so as to permit propulsion of the watercraft in either forward or reverse directions.\nBecause of the fact that the outboard drives, either outboard motors or inboard-outboard drives, may be used either singly or in pairs, it is very desirable if the same basic construction can be utilized for both single and twin installations. This presents problems in connection with twin installations since, as aforenoted, the drives should rotate in opposite directions in such applications.\nIt is, therefore, a principal object of this invention to provide an improved and simplified arrangement for permitting reverse rotation of a marine outboard drive.\nIt is further object of this invention to provide a marine outboard drive that lends itself to ease in reversing the direction of rotation without necessitating major changes to the overall construction.\nIn conjunction with most conventional outboard drives, they are designed so that the input shaft rotates in a constant direction and the driving thrust on the unit always apply in the same direction. However, when the drive is designed so as to be driven in reverse directions for facilitating application with twin drives, the previously proposed constructions have not been completely satisfactory.\nIt is, therefore, a still further object of this invention to provide a marine outboard drive which can easily be rotated in either of two selected directions and wherein the mechanism is designed so as to take loadings regardless of the direction of drive."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic detector for detecting an external magnetic field using a magnetoresistance effect element layer or the like. More specifically, the present invention relates to a magnetic detector that enables highly accurate magnetic detection by adjusting an electrical resistance value, and a method for manufacturing the magnetic detector.\n2. Description of the Related Art\nA magnetic detector is known. In this magnetic detector, a magnetoresistance effect element layer having an electrical resistance which changes in response to an external magnetic field is connected in series to a resistance element layer having an electrical resistance which does not change in response to an external magnetic field, a power supply voltage is applied to the magnetoresistance effect element layer and the resistance element layer, and a potential at a midpoint between the magnetoresistance effect element layer and the resistance element layer is detected, whereby an external magnetic field is detected. A magnetoresistance effect element layer having an electrical resistance which changes in response to an external magnetic field may be used as the foregoing resistance element layer.\nIn this type of magnetic detector, the electrical resistances of the magnetoresistance effect element layer and the resistance element layer are adjusted so that the potential at the midpoint is half the power supply voltage when there is no external magnetic field. When an external magnetic field is generated, the electrical resistance of the magnetoresistance effect element layer changes, and thereby the potential at the midpoint changes. In this way, an external magnetic field is detected by measuring a change in potential at the midpoint.\nMethods for adjusting the electrical resistances of the magnetoresistance effect element layer and the resistance element layer are disclosed in Japanese Unexamined Patent Application Publications No. 2008-46076 and No. 2008-58183.\nIn the method for adjusting electrical resistance disclosed in Japanese Unexamined Patent Application Publication No. 2008-46076, as illustrated in FIG. 21, a plurality of current bypasses, i.e., a first current bypass 601, a second current bypass 602, and a third current bypass 603, are provided to a resistance element layer 703 having an electrical resistance which does not change in response to an external magnetic field, in a stepwise form along the longitudinal direction of the element (X2-X1 direction).\nAlso, a protruding region 604 that protrudes in a direction away from the resistance element layer 703 is provided in each current bypass.\nWhen a difference in resistance value between a magnetoresistance effect element layer 701 having an electrical resistance which changes in response to an external magnetic field and the resistance element layer 703 is outside an allowable range, current paths in the protruding regions 604 are cut off in order from the protruding region 604 that protrudes the most viewed from the resistance element layer 703, and thereby the difference in resistance value is adjusted to be within the allowable range.\nIn this method, a plurality of protruding current bypasses are provided to the resistance element layer 703, which inhibits a reduction of the size of a magnetic detector.\nIn the method for adjusting electrical resistance disclosed in Japanese Unexamined Patent Application Publication No. 2008-58183, as illustrated in FIG. 22, magnetoresistance effect element layers 701 and 702 having an electrical resistance which changes in response to an external magnetic field and resistance element layers 703 and 704 having an electrical resistance which does not change in response to an external magnetic field are connected in series, respectively. Also, there are provided output conductive layers 705 and 706 that are connected to midpoints between the magnetoresistance effect element layers 701 and 702 and the resistance element layers 703 and 704, respectively, and that output potentials of the midpoints.\nEach of the output conductive layers 705 and 706 extends in one longitudinal direction of the magnetoresistance effect element layers 701 and 702 and the resistance element layers 703 and 704. Each of the output conductor layers 705 and 706 is electrically connected to one of the magnetoresistance effect element layers 701 and 702 and the resistance element layers 703 and 704 via a conductive connection layer 707 or 708. In this way, the positions of the connection layers 707 and 708 are changed, so that the electrical resistance value of one of the magnetoresistance effect element layers 701 and 702 and the resistance element layers 703 and 704 is adjusted.\nIn the adjustment of changing the positions of the connection layers 707 and 708, a step-and-repeat exposure system (stepper), which performs exposure using a reticle by moving a substrate using a step-and-repeat method, is typically used for forming the pattern of the connection layers 707 and 708. Thus, the pattern of the connection layers 707 and 708 is formed through a plurality of exposure operations. In the plurality of exposure operations, variations in dimension and variations in alignment of the connection layers 707 and 708 occur. Accordingly, highly accurate adjustment of electrical resistance cannot be performed."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for the metathesis of olefins. More particularly the invention relates to a process wherein the catalyst is part of a distillation structure and the products are simultaneously separated from the reactants and each other by fractional distillation during the reaction.\n2. Related Art\nBroadly metathesis has been defined as a chemical reaction in which an element or radical in one compound changes places with another element or radical in another compound. See The Van Nostrand Chemist\"\"s Dictionary, D. Van Nostrand Company, Inc., 1953, page 463. More specifically olefin metathesis can be defined as the redistribution of alkylidene moieties to give a mixture of olefins. In effect, this reaction takes place via cleavage of the olefin double bond. Generally the reactions of olefinic molecules in the presence of metal-containing catalysts to produce other olefinic molecules are known in the art as xe2x80x9cdisproportionationxe2x80x9d, xe2x80x9cdismutationxe2x80x9d or xe2x80x9cmetathesisxe2x80x9d reactions.\nThe metathesis reactions are of considerable interest because of the versatility of the reaction and the numerous olefinic hydrocarbons available from petrochemical sources which are suitable for use in the reaction to yield useful products. One such reaction is the metathesis of propylene with itself to produce n-butene and ethylene. See, for example, U.S. Pat. No. 4,046,832. The reverse reaction is the metathesis of ethylene with n-butene to produce propylene and is disclosed in U.S. Pat. No. 5,026,936. Another use of the metathesis reaction is to produce 2-methyl-2-butene from the reaction of 2-methyl-1-propene and 2-butene. See U.S. Pat. No. 3,702,827.\nFinally, the simultaneous disproportionation of olefins and fractional distillation is reported in U.S. Pat. No. 4,709,115. Therein the disproportionation of butene with itself to produce ethylene or propylene and hexene or pentenes is disclosed.\nCatalysts that are known to catalyze the metathesis include the oxides of tungsten, rhenium and cobalt/molybdenum.\nBriefly the present invention relates to metathesis reaction carried out a in distillation column reactor, that is, reaction and fractional distillation of the reactants and products are carried out concurrently in the distillation column reactor wherein the catalyst may be in the form to act as a distillation structure or part of a distillation structure or alternatively the catalyst may be located in beds or zones preferably located within the distillation column reactor.\nOne embodiment of the present invention comprises the production of propylene from the reaction of 2-butene with ethylene in a distillation column reactor. Preferably the catalyst is supplied in the form to act as a distillation structure or part of a distillation structure and loaded into the upper portion of a distillation column reactor. The 2-butene is fed above the bed and the ethylene is fed below the bed. Product propylene is taken as overhead and any heavies produced are removed as bottoms.\nIn a second embodiment the present invention relates to a process for the production of higher molecular weight olefins useful in the manufacture of detergents. More particularly the invention relates to the metathesis of olefins having higher molecular weight than that desired with lower molecular olefins to produce the desired molecular weight olefins.\nIn another embodiment the present invention comprises the metathesis of isobutylene with 2-butene to produce 2-methyl-2-butene and propylene. The catalyst is loaded into the upper portion of a distillation column reactor. The isobutylene is fed below the bed. Product 2-methyl-2-butene is taken as bottoms while propylene is removed as overheads.\nAnother embodiment of the present invention comprises the metathesis of isobutylene with itself to produce 2,3-dimethyl-2-butene (tetramethylethylene or TME) and ethylene. The catalyst is loaded into the upper portion of a distillation column reactor. The isobutylene is fed below the bed. Product TME is taken as bottoms along with any C8= or C12=oligomerization product. Ethylene is removed as overheads. The TME is separated from the heavy oligomers by fractional distillation. The C8=oligomers are then separated from the C12=oligomers by fractional distillation and fed to a single pass down flow reactor with ethylene for metathesis to 3,3-dimethyl-1-butene (neohexene)."}
{"text": "Perfluorocyclobutane (C-318, b.p. -6.degree. C.) is a valuable material which can be used as a propellant, etch gas and fire extinguishant. This compound is typically made by cyclodimerization of tetrafluoroethene (TFE) or is recovered as a by-product from the manufacture of TFE. As a result C-318 may be contaminated with unsaturates such as E- and Z-perfluoro-2-butene (i.e., CF.sub.3 CF.dbd.CFCF.sub.3 or FC-1318my) and perfluoroisobutene (i.e., (CF.sub.3).sub.2 C.dbd.CF.sub.2 or PFIB or FC-1318mmt). These compounds, especially the latter, are highly toxic and must be removed prior to commercial use. FC-1318's are difficult to separate from C-318 by distillation as the boiling points are similar.\nU.S. Pat. No. 5,001,287 discloses a process for treating an impure mixture consisting essentially of at least one olefinic impurity and at least one saturated halocarbon by contacting the mixture with a source of hydrogen in the presence of a hydrogenation catalyst. The saturated halocarbon includes perfluorocyclobutane and the catalyst includes Group VIII metals and rhenium. The olefinic impurity is converted to a hydrogenated form. For example, CF.sub.3 CH.sub.2 F containing CF.sub.2 .dbd.CHCl impurity may be treated to produce CF.sub.3 CH.sub.2 F containing CHF.sub.2 CH.sub.2 Cl.\nThere is a need for alternative methods of purification. In particular, in the case of purifying certain perfluorocycloalkanes, there is a need for a process for obtaining products of reaction of the olefinic impurity which are not simply products of adding hydrogen at the point of unsaturation."}
{"text": "1. Field of the Invention\nThe present invention relates, in general, to secure communications, and, more particularly, to secure data communications with a mobile computer over an insecure network.\n2. Relevant Background\nA typical computing environment includes a secure network, such as a local area network (LAN) or wide area network (WAN) that can only be accessed by computers that are authorized by the network administrator to have access. These networks are non-public and so security can be readily controlled with conventional password management techniques. Mobile users can access the network through, for example, dial-up connections through server or gateway that verifies the user\"\"s identity and access privileges.\nAn important use of the Internet and other public data communication networks is the ability to exchange data between mobile computers and an organization\"\"s secure internal network. However, the public network is not secure. An internal secure network uses a gateway machine or xe2x80x9cfirewallxe2x80x9d to couple the internal network to the external insecure network. A firewall is a hardware and/or software system designed to prevent unauthorized access to or from a private network. A firewall examines all packets entering and exiting the private network and blocks those that fail to meet specified security criteria. In an Internet environment, the gateway performs security operations on the IP layer by using, for example SunScreen(trademark) SKIP, (SunScreen is a trademark of Sun Microsystems, Inc.). SKIP is a public key certificate-based key-management scheme which provides key-management for Internet protocols. Data communications using a secure gateway in this manner are referred to as xe2x80x9csecure IPxe2x80x9d.\nAll external hosts must be able to communicate with the internal network using secure IP at any time, but must also be allowed to reach the internal network while transmitting in the clear. This is useful if some services on the internal network must be accessible by the general public (e.g., web server or software download access) and by privileged users such as employees which may have additional rights on those services, e.g., downloading proprietary information. Because of this, a gateway device cannot always provide authorization control simply by filtering out transmissions received in the clear.\nPrior secure IP systems provide authorization control using access control lists (ACLs) that list each IP network address (or other unique network identifier) that is authorized to access a particular resource on the internal network. In general, a gateway can place a static IP address on its ACL and authorize communication from that address to access services on the internal secure network. While this system addresses some problems related to access control, it does not authenticate that the received data packet truly originated from a particular machine.\nA particular difficulty arises in that hosts coupled to the external network may be both regular xe2x80x9cstaticxe2x80x9d Internet nodes (i.e., having a permanently assigned IP address) or mobile nodes (i.e., nodes having a dynamically assigned IP address). It is also possible for a host with a static address to be in secure mode at some time, and be in a clear mode at some other time (e.g., the host running Windows((trademark)) and Unix at different times). Moreover, two mobile hosts with different security properties may appear under the same dynamically assigned IP address at different times. In these instances merely relying on authorization based on the incoming packet\"\"s IP address is insufficient. The gateway machine must be able to authenticate or verify that data received from a remote system truly originated from that system. This situation must be correctly handled by the gateway to prevent, for example, hijacking of TCP connections.\nFor example, when an outside machine using secureIP disconnects from the Internet, thereby relinquishing its IP address, it can be replaced by a second machine transmitting in the clear that has been assigned the first machine\"\"s IP address. From the secure network\"\"s perspective, the incoming TCP packets may have come from either a second machine using the first machine\"\"s IP address, or from the first machine that is now sending in the clear. The second machine will not be able to break the secureIP security, but it may be able to send data in the clear that will reach the internal network. Desirably, the gateway must detect the difference between these two situations, and hinder the second machine\"\"s attempts to send packets on behalf of the old machine. At the same time, the gateway must not allow the fallback to clear text to be abused by an enemy to force all communication to go on in the clear. However, the incoming IP packets do not identify any machine-specific information that would enable the gateway to distinguish between the first machine and the second machine using the same IP address.\nMany proposed approaches to mobile user security require the mobile user to specially configure the security software on the mobile machine. However, this makes the security software more difficult to install and use which is undesirable. To encourage widespread use of secureIP on a variety of machines, it is desirable that the software devices install out of the box, without significant effort to specially configure the software.\nPrior solutions, including SKIP and similar IP security protocols, offer support for mobile hosts by either assigning them a permanent ID (called a master key ID or MKID in SKIP) that is stored in the mobile machine and is transferred with every IP packet. Alternatively, a new security association may be established each time a new mobile IP address is acquired. Although these solutions prevent an intruder with a hijacked IP address from reading encrypted packets, they do not solve the problem of address hijacking so long as the gateway allows the mobile host to send data in the clear. In these cases, the intruder may set the MKID field to zero to force communication in the clear while the security association is maintained by the gateway.\nMoreover, this approach does not allow machines on the internal network to find out whether the incoming link is secure. The gateway holds the list of authorized addresses and performs the encryption/decryption functions. This information is not transmitted or shared with the internal network devices. Hence, the internal network machines cannot tell from examining the header of a received packet whether the packet was from a secureIP link or received in the clear. It would be useful for the internal devices to be aware of this information so that they could take intelligent action in response to receiving a packet with unexpected security properties.\nAnother approach uses xe2x80x9cfirewallsxe2x80x9d which give the capability to do address translation for topology hiding. This hinders non-authorized user\"\"s efforts to find out about the structure and potentially vulnerable points of the internal network. Although this approach makes address hijacking less effective, it does not prevent its occurrence. Another solution relies on control messages transmitted from mobile hosts to establish IP tunnels. These tunnels provide a mechanism needed to redirect data addressed to the mobile host to a dynamically assigned IP address. Tunnels hinder address hijacking by encrypting packet header information as well as the packet payload, but are difficult to set up and require complex security management mechanisms.\nThe Internet Engineering Task Force (IETF) working groups for mobile IP have focused on one potential solution for the support of mobile hosts in the current Internet structure. For this, mobile hosts get assigned a xe2x80x9chome IP addressxe2x80x9d, and a temporary routing address that is used to address traffic. In the gateway from the mobile network to the traditional Internet, address translation and rerouting may be performed, such that the mobile node appears to be reachable on its home address at all times. This approach can result in a security risk if a request message was sent by a hosts that had hijacked the dynamic IP address without cryptographically verifying the authenticity of such messages. In order to avoid this risk, all request messages transmitted by a mobile host to the secure network must be authenticated using a message authentication code such as, for example, the keyed-MD5 algorithm.\nA need exists for a security method and system that support mobile hosts in a public network that solves the security risks created by dynamic IP address assignment to prevent an external machine from impersonating a secured machine, allow internal machines to detect whether the outside machine is coming in using a secured connection, and enable the system to be easily configured and used such that it can bootstrap with little or no user intervention. Desirably, the security method and system can be implemented without access control lists, timers, or other complex security management systems such that it is compatible with load balancing mechanisms.\nBriefly stated, the present invention involves a method for secure data communication between an inside network with a mobile machine in which a data packet is received from the mobile machine having a particular network address. A pool of secure addresses is established and a data structure is created to hold address translation associations. Each association is between a particular network address and a particular one of the secure addresses. If the received data packet is a secure data packet an association between the received data packet\"\"s network address and a secure address in the data structure is identified and the data packet\"\"s network address is translated to the associated secure address before forwarding the data packet on to higher network protocol layers. When the received data packet is not secure it is passed it on without address translation to the higher network protocol layers.\nWhen packets are received by the gateway from the inside network, and are addressed to a secure address, then the secure address is replaced by the corresponding network address and the packet is encrypted and authenticated. As used herein, the term xe2x80x9csecuring a packetxe2x80x9d means authentication and/or encryptionxe2x80x94and not necessarily encryption only. In this manner, bidirectional secure communications are supported."}
{"text": "1. Field of the Invention\nThe present invention relates to a discrete track magnetic recording medium having one surface or both surfaces of a disk-shaped substrate thereof formed with an annular data-recording area having a plurality of tracks separated from each other by concentric separating grooves or a spiral separating groove.\n2. Description of the Related Art\nToday, as a recording medium on which high density recording can be performed, there has been developed a discrete track magnetic recording medium (e.g. a magnetic recording medium disclosed in Japanese Laid-Open Patent Publication No.H04-310621. This magnetic recording medium has a disk-shaped substrate formed with a data-recording area having a plurality of data-recording tracks (hereinafter simply referred to as “tracks”) magnetically separated from each other by a plurality of concentric grooves. This kind of discrete track magnetic recording medium (hereinafter also referred to as the “discrete track medium”) is mounted e.g. on an HDD (Hard Disk Drive), and various record data are magnetically recorded thereon via a magnetic head for use in recording or reproduction (hereinafter also referred to as the “magnetic head”). In this case, in the discrete track medium, grooves are provided between tracks to thereby prevent undesired increase in effective recording track width, due to the leakage field from the magnetic head. Further, the discrete track medium has the advantage of reduced occurrences of recording of record data on adjacent tracks and reproduction of record data recorded on adjacent tracks (cross-talk). This makes it possible to reduce the pitch at which tracks are to be formed (hereinafter referred to as the “track pitch”), whereby record data can be recorded with high density.\nOn the other hand, in an HDD having this kind of discrete track medium (e.g. a discrete track medium 51 shown in FIG. 7) mounted thereon, during recording or reproducing of record data, a motor causes the discrete track medium 51 to be rotated at high speed. Further, as shown in FIG. 8, a servo control mechanism, for example, moves (performs servo control of) a magnetic head 71 along the radius of the discrete track medium 51 such that the central portion of the magnetic head 71 is placed on any one of tracks 62 in a data-recording area 53 formed on a disk-shaped substrate 52. In this case, an air current is generated between the magnetic head 71 and the discrete track medium 51 (data-recording area 53) by the high-speed rotation of the discrete track medium 51, which causes the magnetic head 71 to fly by a lift generated by the air current. As a result, it is possible to move (rotate) the discrete track medium 51 and the magnetic head 71 such that they are prevented from contacting each other.\nHowever, from the study of the proposed discrete track medium 51, the present inventors found out the following points for improvement: In the discrete track medium 51 of the above-mentioned kind, as shown in FIG. 7, an outer area 54 and an inner area 55, neither of which has any grooves 61 (see FIG. 8) formed therein, exist outward and inward of the data-recording area 53 on the disk-shaped substrate 52, and hence inconveniences are caused by the areas 54 and 55. More specifically, as shown in FIG. 9, when data is recorded on or reproduced from a track 62A at the outermost location in the data-recording area 53, it is required to place the central portion of the magnetic head 71 on the track 62A, so that an outer portion of the magnetic head 71, outward of the central portion of the same, protrudes into a region above the outer area 54.\nIn this case, as shown in FIG. 9, the volume of a space defined between the outer portion of the magnetic head 71 and the outer area 54 is smaller than the volume of a space defined between the inner portion of the magnetic head 71 and the data-recording area 53, by an amount corresponding to grooves 61 that are not formed in the outer area 54. Now, the aforementioned lift applied to the magnetic head 71 tends to become larger, as the amount of air between the magnetic head 71 and the discrete track medium 51 is smaller, that is, the volume of a space defined therebetween is smaller. Therefore, as shown in FIG. 9, in a state where the magnetic head 71 is parallel to the surface of the discrete track medium 51, a lift F5 applied to the outer portion of the magnetic head 71 becomes larger than a lift F6 applied to the inner portion of the magnetic head 71. As a result, as shown in FIG. 10, the magnetic head 71 is inclined until the volume of the space defined between the outer portion of the magnetic head 71 and the outer area 54 becomes equal to the volume of the space defined between the inner portion of the magnetic head 71 and the data-recording area 53 (i.e. until the lift F5 becomes equal to the lift F6). Therefore, as shown in FIG. 10, the discrete track medium 51 suffers from the problem that the magnetic head 71 can contact the surface of the discrete track medium 51 (data-recording area 53) due to the thus caused inclination of the magnetic head 71."}
{"text": "Short-run print jobs, e.g., business cards, letterheads, sell sheets, invitations, announcements, folders, brochures, and marketing materials, are generally printed by commercial printers using relatively small, low cost printing equipment. Because of the set-up time involved in changing from one print job to the next, and the relatively low volumes printed (often less than 1000 units/order), the printing cost is typically relatively high, e.g., $20–50 per thousand square inches (“MSI”). In some cases, several print jobs are manually “ganged” together (consolidated or aggregated) onto a single master, in an attempt to reduce the average set-up time per order. Another strategy for controlling cost, employed by printers of products such as invitations, office stationery, and address labels, is to offer customers a limited selection of papers, formats and colors from which to choose.\nPrinting costs per MSI are much lower for high-volume high-quality full-color publishing and packaging print jobs, e.g., food labels, consumer good packaging, magazines, catalogues and high volume marketing materials. Publishing and packaging printing is generally done using large, expensive offset printing presses (either web press or sheet feeding of large-format paper stock) in a highly automated large-volume manufacturing environment. Because these presses have high set-up and amortization costs, their use has been focused on long print runs that are typical in the packaging and publishing segments of the printing market.\nAttempts have been made to reduce the high cost of short-run printing. Set-up costs may be reduced by using rapid changeover production machinery, digital technologies, thermographic printing, or single-color offset printing. Typically, these techniques assume that each print job is to be processed as a discrete production run subject to economies of scale based on the quantity of that print job.\nAnother approach has been to preprint high volumes of a standard base product (e.g., invitation “blanks” bearing high quality color graphics) using high quality offset printing, and then to overprint variable, custom text (e.g., the text of the invitation) for each order, typically using simpler printing processes and conventional short run printing methods.\nYet another approach has been to reduce the cost of setting up a print job by letting the customer, or an intermediary other than the printer, be responsible for the layout, sales and administration aspects of the customer's order. For example, some companies, such as Hallmark, have provided WYSIWYG (“what you see is what you get”) terminals at which a customer can view a WYSIWYG display of the item to be printed, and then upload information regarding the print job to a local or remote printing site. Another example of this approach is desktop publishing software, which allows a customer to design a print job on-screen.\nComputers have been used to reduce cost and improve efficiency of printing processes, e.g., to make the process of page layout, proofing, approvals and transmission to the printing floor more efficient. For example, in the newspaper and printing industries, on-the-fly page markups have been sent directly to the production floor using digital workflow technology. Prepress software and equipment that automates workflow is also used by printers and graphics professionals. Recently, Internet companies such as Noosh and Impresse have been providing services that improve the efficiency of buyer-seller transactions involving printing, e.g., by giving users of their websites the ability to “connect” with a wide variety of print vendors, from short-run demand printers to long-run offset printers."}
{"text": "The present invention generally relates to semiconductor devices (here, the invention chiefly relates to semiconductor memory devices, and a semiconductor memory device will hereinafter be described), and particularly to the redundancy technique for relieving defects by replacing defective memory cells by spare memory cells.\nHigh-density integration of semiconductor memory has been highly advanced up to the mass production of dynamic random access memories (DRAM) of 256 mega bits. This high-density integration advancement involves extreme reduction of element size and great increase of the number of elements, thus causing a problem of low yield due to defects. As a countermeasure against this problem, the so-called redundancy technique is known, in which the defective memory cells are replaced, or repaired by redundant memory cells as the spare, or back-up memory cells that are previously provided on a memory chip. Efforts to enhance the efficiency of the defect repair have been made in this technical field. An example of the defect repair technology for DRAM is disclosed in JP-A-2-192100 (laid open Jul. 27, 1990) in which the decision for column-side repair performed according to a row address, and column selection lines are replaced by redundant column selection lines so that block repair can be made. This method is powerful because a large number of defective memory cells can be replaced by a small number of redundant column selection lines.\nFIG. 2 is a schematic diagram of a conventional logic construction for block redundancy. Here, a memory cell group of two regions including defects is replaced by a redundant memory cell group. A memory cell array NMCA and a redundant cell array RMCA are provided and controlled by a repair decision circuit YRC. The memory cell array NMCA has memory cells provided at the intersections of N word lines WLs and M data lines DLs, and the memory cells are selected by a row decoder XDEC and a column decoder YDEC. The redundant cell array RMCA has redundant memory cells provided at the intersections of N word lines WLs and P data lines RDLs, and the redundant memory cells are selected by the row decoder XDEC and a redundant column decoder RYD. The row decoder XDEC decodes a row address AX of n bits and selectively drives one of the 2n, or N word lines. The column decoder YDEC decodes a column address AY of m bits and selects one of the 2m, or M data lines DLs. The redundant column decoder RYD decodes p bits of the column address AY, and selects one of the 2p, or P redundant data lines RDLs. A repair decision result RYH from the repair decision circuit YRC controls the column decoder YDEC and the redundant column decoder RYD. If the repair decision result RYH is xe2x80x980xe2x80x99, the column decoder YDEC is activated to select memory cells within the memory cell array NMCA. If the repair decision result RYH is xe2x80x981xe2x80x99, the redundant column decoder RYD is activated to select redundant memory cells within the redundant cell array RMCA. Thus, the memory cell group of defects DF1, DF2 can be replaced by a redundant memory cell group. A unit of memory cells to be replaced is the area selected by Q word lines and P data lines.\nThe repair decision circuit YRC is formed by two row address comparators AXC, two column address comparators AYC, two dual-input AND gates AND2, and an dual-input OR gate OR2. A pair of one row address comparator AXC and one column address comparator AYC stores one region to be repaired, or replaced. Each row address comparator AXC includes address storage means for storing a repair address of (nxe2x88x92q) bits, and compares it with the (nxe2x88x92q) bits of the row address AX. Each column address comparator AYC includes address storage means for storing a repair address of (mxe2x88x92p) bits, and compares it with the (mxe2x88x92p) bits of the column address AY. The dual-input AND gates AND2 take logic products of coincidence decision results XHC1, XHC2 from the row address comparators AXC and coincidence decision results YH1, YH2 from the column address comparators AYC to produce decision results HC1, HC2 for the two, first and second replacements. The dual-input OR gate OR2 takes a logical sum of these decision results to produce the repair decision result RYH. Since the repair decision circuit is constructed as above, defects at separate column addresses can be repaired according to the row addresses, or replaced by redundant memory cells on the same redundant data line.\nIn the column block redundancy shown in FIG. 2, the row addresses in the first replacement must be different from those in the second replacement. In other words, the repair row address stored in one of the two row address comparators AXC must be different from that of the other. If the same row address were stored in the two comparators, the replacing regions RPD would be one region, or the replaced regions would compete with each other for acquiring the one region irrespective of whether the column addresses of the replaced regions RPO are different or not. Therefore, even though two repair addresses can be stored, it is impossible to repair the case in which two defects occur in different-column-address regions but in the same-row-address regions each of which is selected by Q word lines and P data lines. In order to reduce the probability of that case in which both defects cannot be relieved because the replaced regions RPO compete with each other about taking one replacing region RPD, it can be considered to decrease the number, Q of word lines that are one replacement unit. However, if the number Q is decreased, it is necessary to increase the number of bits, (nxe2x88x92q) of the row address that the row address comparator AXC compares with, so that the circuit scale of the row address comparator becomes large.\nAccordingly, it is desired to contrive a method of effectively repairing a plurality of defects at the same time. That is, it is an object of the invention to provide a semiconductor memory device having a redundancy circuit capable of effectively repairing defects by use of small-circuit-scale address comparators that compare with a smaller number of bits, and by controlling the replacement operation so that the competition between the replaced regions can be avoided.\nAccording to one aspect of the invention, there is provided a semiconductor memory device having a plurality of word lines, a plurality of bit lines arranged to intersect the plurality of word lines, a large number of memory cells arranged at necessary intersections between the plurality of word lines and the plurality of bit lines, a plurality of spare bit lines arranged to intersect the plurality of word lines, a plurality of spare memory cells arranged at necessary intersections between the plurality of word lines and the plurality of spare bit lines, and a redundancy circuit for replacing memory cell groups, including defects, of the large number of memory cells by spare memory cell groups of the spare memory cells, wherein the redundancy circuit has functions to control a first replacement to be made by a first replacing unit, and a second replacement to be made by a second replacing unit that is smaller than the first replacing unit, and to give the second replacement priority when the first and second replacements compete with each other about taking the replacing spare memory cell groups.\nIn other words, the redundancy circuit controls the first replacement to be made by the first replacing unit and the second replacement to be made by the second replacing unit that is smaller than the first replacing unit, and it includes a first decision circuit for deciding about at least a first part of an address provided to select the large number of memory cells, a second address decision circuit for deciding about a second part of the address, and a third address decision circuit for deciding about at least a third part of the address except the second part, whereby when the second address decision circuit produces a miss, the second replacement is not performed but the first replacement is made according to the output from the first address decision circuit, and when the second address decision circuit produces a hit, the first replacement is not performed but the second replacement is made according to the output from the third address decision circuit."}
{"text": "(1) Field of the Invention\nThe present invention relates to image editing apparatuses and, in particular, to an image editing apparatus capable of editing an image.\n(2) Description of the Related Art\nIt is known that there are image editing apparatuses and image editing programs which can easily edit materials such as obtained still image data and video data. For example, Patent Reference 1 (Japanese Unexamined Patent Application Publication No. 2006-157197) discloses a photo movie creating apparatus for creating a photo movie made of still images with dynamic visual effects added thereto, and gives the images various stories."}
{"text": "1. Field of the Invention\nThe present invention relates to a vaccine for protection of mammals from disease caused by Neospora caninum. More specifically, the invention relates to safe and immunogenically effective vaccines for protection of bovines and canines from abortion caused by Neospora caninum. \n2. Brief Description of the Prior Art\nNeospora caninum was first reported by Dubey et al (JAVMA, Vol. 192, No. 9, May 1, 1988) as a Toxoplasmosis-like illness affecting dogs. Neospora caninum was found to be structurally distinct from Toxoplasma gondii and did not react with anti-T. gondii antiserum in an immunoperoxidase test. Dubey et al described the major lesions associated with the organism as meningoencephalomyelitis and myositis. Within the past few years, neosporosis has become recognized as a major reproductive disease in cattle (Anderson et al., 1994, Food Animal Practice, 10: 439-461) with cases reported in North and South America, Europe, Africa, the Pacific-rim countries as well as in the United States. The major clinical manifestation of bovine neosporosis is fetal abortion, with focal nonsuppurative necrotizing encephalitis, nonsuppurative myocarditis, and myositis in the fetus (Anderson et al., 1991, Journal of the American Veterinary Medical Association, 198: 241-244). According to Anderson et al., 1997 (Journal of the Veterinary Medical Association, 210: 1169-1172), retrospective studies of cattle in California indicate that neosporosis has been endemic since at least 1985. These authors state that 18 to 19% of all aborted bovine fetuses submitted to the California Veterinary Diagnostic Laboratory System have Neospora sp infection. In a prospective survey of selected dairies in California, the number of abortions attributed to Neospora sp infections was even greater (42.5%).\nHo et al (J. Parasitol., 1997, 83(3)) have recently reported the successful reproduction of bovine abortion and fetal infection by infecting pregnant cows with tachyzoites of Neospora caninum. This publication suggests that there may be a correlation between serological titer as measured by indirect fluorescent antibody (IFA) testing and protection from abortion caused by Neospora caninum in cows. Cows with IFA titers of 320 and 640 did not abort after infection with tachyzoites of this organism.\nAs mentioned previously, neosporosis has also been reported in puppies and in dogs as old as 15 years of age. The percentage of infected dogs that show clinical signs is unknown. In dogs, Neospora caninum can infect any tissue, although it is most commonly found in the central nervous system and spinal nerve roots. The most severe infections are seen in puppies that were infected in utero. These puppies exhibit ascending paralysis. Abortion can be reproduced in experimental infection of pregnant bitches during the early stage of gestation. Sulfonamides, pyrimethamine and clindamycin have been used to treat neosporosis in dogs.\nNeospora caninum can also produce a fatal infection in experimentally inoculated cats. However, the disease has not yet been reported to occur naturally in cats.\nNeosporosis has been observed to cause abortion in sheep and goats but to a lesser extent than is found in cattle. Experimental infection is readily induced in sheep and goats by subcutaneous injection of tachyzoites.\nAlthough neosporosis, especially in cattle, appears to pose an increasingly serious problem and there is certainly a long felt need to solve this problem by protecting mammals using a vaccine, there are no descriptions of vaccines, vaccine development nor suggestions of methods of preparing vaccines to protect cattle and other animals from disease caused by Neospora caninum. "}
{"text": "Integrated circuits include increasing numbers of memory access requestors such as video processing units, graphical processing units, general purpose processors, peripherals, and the like. The amount of information consumed by these memory access requesters as well as their expected throughput is constantly increasing.\nOn the other hand, the power consumption of the integrated circuit should be maintained at a relative low level, especially in integrated circuits that are powered by batteries. Accordingly, an integrated circuit is expected to provide a trade-off between performance (throughput, speed) and power consumption.\nVarious power reductions techniques include dynamic voltage and frequency scaling (DVFS), and power gating. DFVS alters the voltage supply level and the frequency of clock signals that are provided to a circuit while power gating involves shutting down the entire circuit or selected portions of that circuit.\nThe power consumption of various memory access requesters can be reduced by applying DVFS techniques."}
{"text": "1. Field\nThis disclosure relates to audio signal processing.\n2. Background\nBeamforming is a signal processing technique originally used in sensor arrays (e.g., microphone arrays) for directional signal transmission or reception. This spatial selectivity is achieved by using fixed or adaptive receive/transmit beampatterns. Examples of fixed beamformers include the delay-and-sum beamformer (DSB) and the superdirective beamformer, each of which is a special case of the minimum variance distortionless response (MVDR) beamformer.\nDue to the reciprocity principle of acoustics, microphone beamformer theories that are used to create sound pick-up patterns may be applied to speaker arrays instead to achieve sound projection patterns. For example, beamforming theories may be applied to an array of speakers to steer a sound projection to a desired direction in space."}
{"text": "As a result of the development of machines for automatically inserting electronic components into printed circuit boards, it has become desirable to provide mechanisms to automatically supply the board to the component insertion machines. Devices of this type are illustrated in U.S. Pat. Nos. 3,722,062; 3,808,662; 4,209,898 and 4,283,847. These mechanisms all operate on the principle of transferring the board from one end of the component insertion device to the other by means of a continuous conveyor of some type upon which the board is carried. This requires the board to be removed from one storage magazine and transferred to a separate receiving magazine after the component insertion has taken place."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nA vehicle may include an internal combustion engine and/or one or more electric motors that generate torque. A transmission selectively transfers torque to one or more wheels of the vehicle. An engine control module (ECM) controls operation of the engine. The ECM or another control module, such as a motor control module, may control an electric motor. A transmission control module (TCM) controls the transmission. The vehicle may also include one or more other control modules, such as a chassis control module, etc.\nThe control modules may communicate with one another via a network that can be referred to as a car area network (CAN). The control modules may communicate, for example, to share data. One or more of the control modules may make a decision and/or take action based on the shared data.\nThe vehicle also includes a communications module that is connected to the car area network. The communications module may also selectively establish a connection to a remote server. The communications module may update one or more of the control modules based on data from the remote server. The communications module may selectively output one or more vehicle parameters to the remote server."}
{"text": "1. Field of the Invention\nThe present invention relates to a gaming system, more specifically to a gaming system with a plurality of terminals where players can make a selection from a plurality of games or a plurality of images to be displayed on a display, which collects various information regarding the selected games or the selected images and provides a visual display desired by the players for each time of day.\n2. Related Art\nInformation such as a player's behavior and preference is extremely important in game halls having a plurality of terminals to provide games. Generally, game halls provide a plurality of types of games; however, space for installing gaming machines is limited. Therefore, games to be installed in the game hall and the allocation thereof have a great influence on operation of the game hall. A game hall with an unpopular gaming machine may discourage players from visiting.\nIn addition, especially with video games, operating rates may be different between similar games; for example, those with popular characters displayed may have a higher operating rate than those with banal character displayed.\nAs described above, installation of game machines in the game hall is directly linked to operation and profit of the game hall, and thus plays an important part in management of the game hall. In order to solve the abovementioned problems, a game hall simulation system, as disclosed in U.S. Pat. No. 6,007,425, is known. With the game hall simulation system, a game hall manager can simulate the operational status of the game hall and traffic route of players. The manager can decide the location of game machines based upon a result thereof, their own experience and the like.\nHowever, in reality, players of different generations and different preferences visit the game hall, depending on the operating hours of the game hall. In such a case, simulating each case with a simulation system to decide which gaming machines to install may take too much time and can not provide a flexible response.\nEven with terminals that allow players to select a game, less experienced players tend to choose a gaming machine to play by looking at a display thereof and judging that the gaming machine provides a desired game. Therefore, if a stand-by screen of the gaming machine does not show the player's desired game, the players may misunderstand that the desired game is not provided in the game hall and leave the game hall.\nIn order to solve the abovementioned problems, the present invention provides a gaming system that allows players to select a game from a plurality of games, storing information regarding games and images selected in each terminal to a server, and provides a display for a game most selected in the same unit time; thus providing an appropriate gaming machine for players at each time of day."}
{"text": "Conventionally, as disclosed in, for example, Japanese Patent Laid-Open No. 2006-287321 (JP '321) and Japanese Patent Laid-Open No. 2011-087026 (JP '026), a vehicular communication system that establishes a narrow area communication channel, such as a Bluetooth communication channel (Bluetooth: a registered trademark), to a cellular phone that is brought into the vehicle compartment is known.\nFurther, as disclosed in Japanese Patent Laid-Open No. 2003-309865 (JP '865) a wireless communication apparatus establishing a wide area communication channel to a wireless communication access point, or a so-called hot spot, at places such as a coffeehouse, a convenience store, a dealer, which are outside of a vehicle, is also known.\nThe communication standards for establishing a wide area communication channel, such as IEEE802.11b, IEEE802.11g, and IEEE802.11n, use a communication band of 2.4 GHz, which is also used for establishing a Bluetooth communication channel. Therefore, simultaneous use of two types of communication, that is, the wide area communication channel according to IEEE802.11b, IEEE802.11g, or IEEE802.11n and the narrow area communication channel by the Bluetooth standard, may require a work-out scheme, such as a separate/distant positioning of respective communication antennas, a time division communication control or the like.\nHowever, such work-out scheme may lead to a deteriorated throughput of the wide area communication channel, due to the communication operation of the narrow area communication channel. More specifically, the wide area communication channel assumes a large amount of data communication to and from the access point, thereby making it necessary and urging for the communication apparatus, which is capable of establishing both of narrow and wide area communication channels, to prevent the throughput deterioration of the wide area communication channel caused by the communication operation of the narrow area communication channel."}
{"text": "1. Field of the Invention\nThis invention relates to the verification of, based on emulation, system-level digital circuits which have been designed; more particularly, relates to a prototyping system which implements rapidly a designed digital circuit in the same form of actual hardware with \"Reusable Field Programmable Devices\" (or chips) and other discrete devices and verifies the designed digital circuit in a short time, and to a method for operating the prototyping system.\n2. Description of Prior Art\nAs technologies for the design of integrated circuits and the manufacture of semiconductor improve rapidly, the scale and complexity of digital circuit design has increased tremendously, and it has been possible to design and manufacture system-level digital circuits capable of implementing microprocessors, digital signal processors, memory devices, application specific logic circuits, etc. on a single chip. Meanwhile, since it becomes necessary to develop a new high quality product in a very short time in order to survive the vicious competition in the market, high-speed prototyping systems capable of implementing and verifying designed circuits effectively have become increasingly important.\nUp to now, a logic simulator based on software has been used to verify the operation of designed digital circuits; however, the logic simulator is not absolutely reliable since it is based on software modeling in a virtual environment, not on implementing an actual circuit. The capability of simulation software and computer cannot affiliate the size and complexity of digital circuits to be verified. To overcome the shortcomings of the logic simulator and to verify a designed circuit more accurately, it has been developed a technology for verifying a designed circuit based on hardware emulation which is closer to the actual digital circuit and which technology is capable of verifying the designed circuit more accurately (hereinafter referred to as \"Rapid Prototyping\"). The verification of a designed digital circuit by using the \"Rapid Prototyping\" can be executed at high speed since it is executed by an actual hardware implementation. The verification of a designed digital circuit by using the \"Rapid Prototyping\" also allows the design verification to be performed in an environment close to the actual hardware. The verification process for a digital circuit design using this prototyping is such that verifying signals are applied at the input of hardware on where the designed circuit to be verified is implemented, and logic values appearing on many signal lines existing in the circuits including output lines are probed to see if they match the logic values intended by the designer by rigorously probing the signal lines.\nRecently, \"Reusable Field Programmable Devices\" (hereinafter referred to as \"RFPDs\") have been frequently used in industry, since any digital circuit can be implemented in the RFPDs by being simply programmed onto the devices, and further, due to the unlimited programmable number, the devices can be reused. There are several types of RFPDs such as a field programmable gate array, a complex programmable logic device, etc. In general prototyping systems based on emulation, the digital circuit to be verified is implemented in the prototyping engine which is composed by interconnecting several ten to several thousand RFPDs and other discrete devices such as microprocessors, digital signal processors, application specific non-memory devices or memories. The capacity of a prototyping system is determined based on an important factor, which is whether the prototyping system can provide solutions for the following problems. First, probing signal lines is inherently difficult since many signal lines existing within the digital circuits to be verified exist inside RFPDs; this problem becomes more serious when more highly integrated RFPDs are used. Second, in a circuit design using the so called \"Intellectual Property\" (hereinafter referred to as \"IP\") which has begun to be used frequently in system-level designs, it is required that a variety of IPs can be used in convenience by the prototyping system, since the system uses different IPs to execute the verification of a designed circuit effectively. To solve this problem, a method for replacing the existing devices or boards into prototyping engine with devices or boards needed by a user or for newly adding the needed devices or boards should be provided. Third, the usage rate of RFPDs by the prototyping engine is very low when a digital circuit is implemented in the prototyping engine, since the number of pins of RFPDs cannot be increased proportionally regardless of increasing RFPDs' integration degree. Fourth, it is required that low integrated RFPDs used in conventional prototyping engine can be replaced with high integrated RFPDs with ease so as to increase the capability of the prototyping engine.\nAs a conventional prototyping system, for example, U.S. patent publication No. 5,452,231 has been disclosed, and the content of the publication is incorporated in this specification as a reference. FIG. 38 is a schematic diagram of a conventional prototyping system. As shown in FIG. 38, in the conventional system, a conventional prototyping engine 8 is connected to a server computer 40 through an interface module 42. In the prototyping engine 8, a plurality of RFPDs 2 used for implementing logic gates and a plurality of RFPDs 4 which are used for connection are connected to an emulation system bus 6, and the emulation system bus 6 is connected to a system interface bus 26. To form the conventional prototyping engine (8), the RFPDs used in the prototyping engine are divided into two groups, RFPDs used for implementing logic gates and RFPDs used for interconnection of logic gates. The connection between the two groups is implemented in partial crossbar, where there is no direct connection between the RFPDs in the same group. This connection type can be a solution for the timing problems, which occur in logic emulation, but also causes other problems. In the conventional prototyping systems, other devices accept the RFPDs, such as microprocessors, digital signal processors and application specific non-memory devices cannot be used easily in the prototyping engine, and it is not easy to replace the RFPDs used in the prototyping engine with high integrated and high speed RFPDs. Further, in case that signal lines existing in a designed circuit to be verified exist also in certain RFPDs, real time probing is impossible, and too many RFPDs are needed in the emulation since the rate of using the resources in the RFPDs is too low due to the physical limitation of the number of input/output pins of each RFPD."}
{"text": "Databases, such as relational databases, are commonly used for managing large amounts of data, and a database user, such as a developer or end-user, can access, modify or add data within the database by way of a query. For example, in the Structured Query Language (“SQL”), a query may have what is referred to as multi-part identifiers in the conventional “dot notation,” such as “A.B.C.D. ” As the name implies, such identifiers has are composed of multiple parts (e.g., A, B, C and D) separated by “dots,” or periods. Such a multi-part identifier is typically read from right to left, where the rightmost part is a column name, preceded by a table name, schema name and database name. Thus, in the above example, D would be the column name, C the table name, B the schema name and A the database name. Less than the entire path may be specified if the user is currently using the database within one of the specified levels. For example, if the user is currently operating in the table in which the desired information resides, then simply entering D as the query would result in the desired column. If the user is operating in the same database, but in a different schema, then a query of B.C.D would be required.\nAs databases become more advanced, a user is provided with greater flexibility, both in terms of the types of data that may be requested and the manner in which a query may be made to request the data. For example, User Defined Types (“UDTs”) are user-defined, structured data types that may be implemented using the conventional dot notation discussed above. A multi-part identifier of A.B.C may be used to query a database for a “property C of property B of column A” or a “property C of column B of table A.” These are in addition to the “legacy” multi-part identifier meaning discussed above.\nA database that is attempting to resolve the A.B.C multi-part identifier discussed above does not know which of the above-noted interpretations the user is intending to achieve. In most cases, only one of the interpretations will achieve a result (i.e., will “bind successfully”). This is because typically there will only be one interpretation that will bind for every identifier specified in a query. Thus, one interpretation of the A.B.C multi-part identifier may look for column C of table B of schema A. If one or more of the identifiers do not correspond to a value in the database, then the database will assume that this interpretation is not correct. The database may then attempt to bind an interpretation of the A.B.C query that looks for a property C of property B of column A, or for a property C of column B of table A, and so forth. If only one of the interpretations binds successfully, the database may assume that such an interpretation is the user's desired interpretation, and may return a result.\nUnfortunately, in some scenarios more than one interpretation may bind successfully. Using the above A.B.C multi-part identifier example, if the database being queried has both a column C of table B of schema A, and a property C of property B of column A, the database will not know which interpretation was desired by the user. Conventionally, this conflict was resolved by using the “most local” interpretation. In other words, the interpretation that contains the least number of column prefixes would be assumed to be the correct interpretation. In the above example, such an assumption would give the highest precedence to the interpretation of “property C of property B of column A,” followed by “property C of column B of table A,” and then “column C of table B of schema A.” Once this assumption is made, the result is returned to the user.\nSuch an algorithm has a significant shortcoming. Namely, if the assumption returns an incorrect result (e.g., property C instead of the intended column C), the user may not know that the returned result is incorrect. This shortcoming is particularly acute when the returned result is in the same format as the intended result. For example, if the returned and intended results are both numerical values, there may be no way for a user to realize that the returned result is incorrect by merely looking at the result. The incorrect result might then go on to corrupt other data if it is used in calculations or the like. Tracking down the source of such an incorrect result may prove to be extremely difficult for a user because there will be no outward sign that the multi-part identifier resulted in two possible interpretations.\nIn other scenarios, none of the possible interpretations bind successfully. Conventionally, in such a case a generic error message such as “could not find A.B.C” would be issued. Such a message is typically not very helpful to a user because it does not describe the interpretation(s) that were attempted, and what the exact problem with each interpretation was.\nAccordingly, it is desirable to have a mechanism for discerning user intent when interpreting a multi-part identifier such as a UDT. In addition, it is desirable to have a mechanism that provides useful error messages to a user in the event of an interpretation conflict, failure or the like. The present invention satisfies these needs."}
{"text": "Diisocyanates which contain urea groups are generally distinguished by a relatively high melting point. Due to the low vapor pressure thereof and due to the presence of the urea groups, they are preferably used as starting materials for the production of highgrade polyurethane/polyurea plastics. These diisocyanates are usually produced by reacting 2 mols of diisocyanate with 1 mol of water in an inert solvent. It is preferred to use diisocyanates which have two differently-reactive groups, when a definite reaction course is to be achieved. Thus, French Pat. No. 1,103,329 describes the production of bis-(3-isocyanato-4-alkoxy-, -alkyl- or chloro-phenyl) ureas by a reaction of diisocyanate and water in a molar ratio of 2:1, in a solvent, which solvent must not contain reactive hydrogen atoms (for example in the form of hydroxyl groups). Ethyl acetate, diethyl ether or acetone are preferably used as solvent.\nU.S. Pat. No. 3,906,019 describes a process for the production of monomeric di-(isocyanatotolyl)-urea, in which the reaction is carried out in an excess of one of the two non-miscible reactants (toluylene diisocyanate and, preferably, water) at from 20.degree. to 40.degree. C. The urea which is not soluble in either of the starting materials is separated in a yield of from about 60 to 85%, mainly in monomeric form, and it must then be purified. Lewis acids or Lewis bases may be used as catalysts for the reaction described therein, but pyridine is preferred\nAll these processes are carried out in solvents. The diisocyanates must be very soluble and the water which is added must be at least partly miscible in the solvents. The solvent must not exert a polymerizing effect on the isocyanate and must be free from isocyanate-reactive functional groups. The urea diisocyanates are produced as relatively insoluble compounds and are isolated by filtration. For further processing of these urea isocyanates for polyurethane production, it is necessary to convert the urea isocyanates obtained by filtration (and freed from solvents under vacuum) into a finely-divided form by a suitable grinding operation. As a result of the high melting point and the difficult solubility of these isocyanate ureas, inhomogeneous products are often obtained during the reaction.\nAccording to German Auslegeschrift No. 2,902,469, suspensions of (poly)-urea-diisocyanates in isocyanateprepolymers are produced by dissolving 1 mol of aromatic diisocyanate in the NCO-prepolymer and reacting it with from about 0.4 to 0.8 mols of water. Organic diisocyanates which have differently-reactive NCO groups are used to produce the NCO-prepolymers, as well as to produce the (poly)urea diisocyanates.\nNo process has been known in which a polyisocyanate containing urea groups is directly produced by reacting organic polyisocyanates with water in an NCO-reactive organic solvent, for example long-chain hydroxyl polyethers or hydroxy polyesters. According to the prior art, such a process would not seem very promising since the hydroxyl groups of the polyether or polyester generally have approximately the same reactivity with respect to NCO groups as water. Thus, it would be expected during the reaction of a polyisocyanate with water in a polyether or polyester for the NCO groups to react to about the same extent with the water and with the OH-groups of the polyol.\nThe separate production of urea-isocyanates, the isolation and purification thereof and the subsequent dispersion or dissolution of the urea-isocyanates in relatively high molecular weight polyols, such as are used for the synthesis of polyurethanes, is expensive. An object of the present invention is to synthesize urea isocyanates directly in the relatively high molecular weight polyols from aromatic isocyanates which are easily accessible in a simple and economic process, in order to be able to use these urea isocyanate/polyol dispersions or solutions directly for the synthesis of polyurethanes."}
{"text": "The present invention relates generally to a housing for a two spool tape or data cartridge for magnetic recording tape of the type used in computers and the like, and is particularly concerned with an improved fastener mechanism or system for coupling or connecting the opposed lower and upper walls of the housing together.\nTypical magnetic tape cartridges include an outer housing and various internal components including tape, a pair of rotatable spools carrying the tape, and a drive mechanism. The housing usually comprises a metallic baseplate and an opposing cover member of plastic material. Where the baseplate and cover member are rigidly secured together, warpage of the housing can occur as a result of the different rates of thermal expansion of the materials forming the baseplate and cover. This warpage results in some deflection of the tape transport path, which is unacceptable in view of the very high track densities on current magnetic data tapes, requiring a high degree of precision in the tape transport.\nIn U.S. Pat. No. 4,571,789 of Morioka, a fastener mechanism for connecting the baseplate to the cover in a manner allowing for some degree of different relative thermal expansion is described. The mechanism comprises a single rigid fastener between the baseplate and cover which is located on the mid-line of the cartridge, and a series of snap fasteners between the baseplate and cover, one each adjacent each corner of the cartridge. The snap fasteners each have a pair of spring arms separated by a slot which can absorb some of the differential expansion between the baseplate and cover by allowing the arms to move towards and away from each other. However, this arrangement only allows relative motion between the baseplate and cover in one lateral direction, parallel to the front of the baseplate, since the spring arms can only move towards and away from each other in this direction, as permitted by the slot. The snap fasteners are actually quite stiff in the front to rear direction. Additionally, the snap fasteners do not provide a secure connection between the baseplate and cover and thus allow some potential for separation of the baseplate from the cover, especially along the front edge of the cartridge where it contacts the drive mechanism and thus encounters relatively large forces."}
{"text": "The present invention relates to an orthodontic appliance for use in the treatment of dental malocclusion.\nThe newest and the most effective method of orthodontic treatment, the main purpose of which is to treat dental malocclusion, is the technique of firmly fixing a bracket to each tooth and inserting and fixing a U-shaped wire into these brackets in order to correct a malalignment of the teeth by the elastic force of this wire. For example, as shown in FIG. 2, when the upper jaw central incisor (1) protrudes abnormally toward outside as shown by the solid line, it is corrected to the state shown by the broken line, by a method in which a band (2) is applied to each tooth (1) as shown in FIG. 1 and brackets (3) are welded to the outsides of these bands (2), i.e. to their labial surfaces or buccal surfaces, then a curved orthodontic wire (4) is fixed to these brackets (3) to correct teeth malalignment by the elastic force of this wire (4). As examples of this method, there are the ribbon arch technique, the edgewise technique, the Begg technique, the universal technique, etc. Among them, the brackets of the edgewise technique are, as shown in FIG. 3, composed of a flange (5) like a metal net or a plate and the tip (6) is attached to this flange (5). A fixing groove (7) is formed in the tip (6), so that a rectangular wire (4) may be inserted horizontally. Wings (8) are formed on both sides of the groove, and a ligature wire (9) retains the wire 4 in the groove. In the Begg technique, brackets (3) are also composed of the flange (5) and the related tip (6), as shown in FIG. 4, but the fixing groove (7) is formed so that a wire (4) may be inserted vertically from the opposite side of the occlusal plane. In addition, a slot is made for insertion of a lock pin (10) shown by the broken line. This pin is inserted from the gingival side and the protruding lower edge of the inserted lock pin (10) is bent and fixed. The ribbon arch technique and the universal technique similar to the edgewise technique or the Begg technique.\nThe main disadvantage of these conventional methods is that the bracket (3) and wire (4) are visible clearly from the outside of the mouth when the mouth is opened, because the brackets (3) are bonded or welded to the outer surfaces of teeth (1) i.e. their labial or buccal sides. This detracts from the facial beauty of the patient and, accordingly, imposes a considerable psychological burden on the patients.\nIn conventional orthodontic treatment, there are some methods for treating malocclusion from the lingual or palatal side. Those are the lingual arch appliance or the palatal plate appliance. Those methods push several teeth simply from the inside by the auxiliary spring attached to the lingual arch or plate. Therefore, the cases which can be treated by those techniques are limited. Accordingly satisfactory results are not obtained in the treatment of all teeth. Neither is it possible to use the conventional bracket itself by bonding same to the inside surface of the teeth. For example, it is not only extremely difficult to bond the bracket for the edgewise technique or Begg technique to the inside surface of the teeth, but also insertion of a wire into the grooves of brackets from the gingival side, fixation by the ligature wire, and insertion and fixation of lock pins into brackets are extremely difficult or impossible. The same difficulties also exist in the ribbon arch technique or the universal technique."}
{"text": "The mixer bowls of horizontal mixers tend to be quite large and therefore commonly have a bowl tilt feature so that mixed product such as dough can be emptied from the bowl (e.g., into a dough trough alongside the bowl). During tilt the bowl pivots on a horizontal axis relative to an overhead canopy that covers the open top of the bowl during mixing. Standard canopy sealing for single tilt (i.e., to only one side of the canopy) mixer bowls is with a straight compression seal across the front of the canopy (where the front is the tilt side) and a scraper on the rear of the canopy to scrape the downward facing surface of the canopy during the tilt. For two-way tilting (i.e., to either side of the canopy), seal bars that pneumatically are positioned and compressed against line of contact between the bowl and canopy have been used. However, the pneumatic sealing arrangement is complex and can create maintenance issues.\nBowl mounted seals are also known, but provide less than desirable results and may increase the possibility of seal mount components entering the food product.\nIt would be desirable to provide a simpler, but effective system for sealing the bowl relative to the canopy, particularly in the case of two-way tilting."}
{"text": "This invention relates to random sequence generators and in particular to a one-dimensional cellular automaton for generating random sequence signals.\nSequences that seem random are needed for a wide variety of purposes. They are used for unbiased sampling in the Monte Carlo method, and to imitate stochastic natural processes. They are used in implementing randomized algorithms on digital computers, and their unpredictability is used in games of chance. Random sequence generators are used as repeatable noise sources for electronic testing purposes.\nTo generate a random sequence on a digital computer, one starts with a fixed length seed, then iteratively applies some transformation to it, progressively extracting a random sequence which is as long as possible. In general, a sequence is considered \"random\" if no patterns can be recognized in it, no predictions can be made about it, and no simple description of it can be found. However, even if the sequence can be generated by iteration of a definite transformation it can nevertheless seem random if no computations done on it reveal this simple description. The original seed must be transformed in such a complicated way that the computations cannot recover it.\nThe fact that acceptably random sequences can be generated efficiently by digital computers is a consequence of the fact that quite simple transformations, when iterated, can yield extremely complicated behavior. Simple computations are able to produce sequences whose origins can apparently be deduced only by much more complex computations.\nMost current practical random sequence generation computer programs are based on linear congruence relations (of the form x'=ax+b mod n), or linear feedback shift registers of the type discussed in \"Shift Register Sequences\", S. W. Golomb, Holden-Day (1967). The linearity and simplicity of these systems has made complete algebraic analysis possible, and has allowed certain randomness properties to be proved. But these characteristics also lead to efficient algebraic algorithms for predicting the sequences (or deducing their seeds), and limits their degree of randomness.\nThere are many standard mathematical processes which are simple to perform, yet produce sequences so complicated that they seem random. An example is taking square roots of integers. Many physical processes also yield seemingly random behavior and, in some cases, the randomness can be attributed to the effects of an external random input. Thus for example \"analog\" random sequence generators such as noise diodes work by sampling thermal fluctuations associated with a heat bath containing many components. Coin tossings and Roulette wheels produce outcomes that depend sensitively on initial velocities determined by complex systems with many components. It seems however that in all such cases, sequences extracted sufficiently quickly can depend on only a few components of the environment, and must eventually show definite correlations.\nI have found that randomness in many physical systems may not arise from external random inputs, but rather through intrinsic mathematical processes. This invention makes use of such a process employing a one-dimensional cellular automaton."}
{"text": "Many vehicles, such as recreational vehicles, currently have furniture that includes a human support surface (e.g., beds, couches, chairs, etc.) that an occupant can sit and/or lay on to occupy the furniture. However, the furniture is typically affixed to one or more surfaces (e.g., the floor and/or one or more walls) of the vehicle. Where the furniture is affixed to the vehicle, the owner of the vehicle and/or an occupant of the furniture cannot relocate the furniture within the vehicle without having to detach the furniture, move the furniture to the desired position, and reattach the furniture to the vehicle.\nAdditionally, where the furniture is affixed to the vehicle, when a person is occupying the furniture, and the vehicle experiences accelerations as a result of one or more maneuvers of the vehicle, the person may shift position because of changes in their inertia, and may be disturbed as a result. This can be problematic, particularly where an occupant of the furniture is performing any sort of task that requires attention to detail and/or fine motor skills."}
{"text": "Mobile communication systems continue to grow in popularity and have become an integral part of both personal and business communications. Various mobile devices now incorporate Personal Digital Assistant (PDA) features such as calendars, address books, task lists, calculators, memo and writing programs, media players, games, etc. These multi-function devices usually allow electronic mail (email) messages to be sent and received wirelessly, as well as access the Internet via a cellular network and/or a wireless local area network (WLAN), for example.\nCellular devices have radio frequency (RF) processing circuits and receive or transmit radio communications signals typically using modulation schemes. In the typical device, the RF processing circuits may include a modulator, a power amplifier coupled downstream from the modulator, and an antenna coupled downstream from the power amplifier. Depending on the immediate surroundings of the cellular device, the impedance load of the antenna may vary, which can impact antenna performance.\nSome cellular devices include an impedance matching device between the antenna and the power amplifier to compensate for the impedance mismatches. One drawback to this approach is that a phase change is imparted onto the transmitted signal. Depending on the wireless protocol being used, this may cause issues with the receiver device."}
{"text": "1. Technical Field\nThe present disclosure relates to methods of analyzing, screening and/or evaluating enzyme(s), polypeptide(s), mixtures of enzymes, mixtures of polypeptides, and/or mixtures of enzymes and polypeptides of interest. The methods are useful for, inter alia, detecting residual cellulose and evaluating hydrolysis activity of enzymes, polypeptides, or mixtures of interest. The present disclosure further relates to high-throughput assay(s), for example, a high-throughput assay for quantifying cellulolytic activity and/or cellulase activity of enzyme(s) or polypeptide(s) of interest in a pretreated corn stover (PCS) hydrolysis is disclosed.\n2. Description of the Related Art\nCurrent methods of analyzing, evaluating, or screening the hydrolysis activity of an enzyme, polypeptide, or mixtures of interest are problematic and not well designed for high-throughput analysis. For example, one method of analyzing pretreated corn stover (hereinafter referred to as “PCS”) hydrolysis requires a lengthy high pressure liquid chromatography (hereinafter referred to as “HPLC”) analysis for determining the quantity of sugars. Here, cellulose, in the PCS is enzymatically hydrolyzed to glucose, cellobiose, and higher beta-glucans. HPLC is used to measure the glucose and cellobiose. One of skill in the art, knowing the cellulose content of the substrate, can then use this information to calculate the percent conversion of cellulose into sugars. Accordingly, cellulolytic activity of the enzymes in the hydrolysis can be measured. However, the HPLC step is time consuming and laborious. Assays using HPLC are not suitable for high-throughput analysis, and/or a quick analysis of multiple enzymes of interests in a single assay.\nAttempts have been made to improve the HPLC assay using a pipettable substrate loaded in deep well plates; however post-hydrolysis HPLC sugar analysis is still required. The HPLC assay is time-consuming, for example, a 96 well plate takes approximately 19 hours of HPLC time. The attempted improvements are problematic in that they include time-consuming filtration, pipetting, and dilution steps prior to HPLC.\nProblems with the known assays lead to higher research costs, tedious assay formatting, as well as time-consuming enzyme activity evaluation. Accordingly, there is a continuing need for assays and analysis methods having improved accuracy and/or reduced performance time, having excellent accuracy, an excellent reaction rate, and/or excellent cellulose conversion, especially where high-throughput analysis is desirable."}
{"text": "Caps are used to prevent fluid leakage from a container filled with fluid. For fluid containers that are dispensed through various dispense apparatuses, the caps are designed to enable drawing the fluid out of the fluid container while the fluid container is connected to a dispensing apparatus. Some caps reseal once the fluid container is disconnected from the dispensing apparatus to prevent spillage. Furthermore, the caps are used to prevent the fluid inside the fluid container from being contaminated with dirt or other substances that can reduce the fluid quality.\nKnown dispensing apparatuses rely on gravity to pump the fluid from the fluid container through the cooler to the faucet of the dispensing apparatus. Thus the fluid container must be connected to the top of the dispensing apparatus. Some fluid containers comprise a sealing component that is removed before the fluid container is connected to the dispensing apparatus. Some of these sealing apparatuses are removed when the fluid container is connected to the dispensing apparatus. When the sealing apparatus is removed by a rigid tube, the fluid container is placed on the top of the dispensing apparatus and the force of gravity pushes the fluid container down on the rigid tube. When the rigid tube comes in contact with the sealing component, the rigid tube pushes out or breaks the sealing element so that fluid may flow into the dispensing apparatus.\nSome capping apparatuses, such as the commercially available cap of Scholle provide for a silicon cap attached to a retaining ring. The silicon cap is flexible and comprises an aperture via which a pin enters to allow passage of fluid via the silicon cap."}
{"text": "Full coverage of a second generation (2G) or third generation (3G) network such as Global System for Mobile Communications (GSM) or Wideband Code Division Multiple Access (WCDMA) has been basically implemented.\nWith development of Long Term Evolution (Long Term Evolution, LTE) network technologies, LTE networks have covered some urban areas and traffic hotspot areas. In this way, in current communications networks, the LTE network coexists with the 2G or 3G network.\nIn a call process, user equipment (UE) may enter a 2G or 3G network from an LTE network. Because the LTE network has a different bearer mechanism from the 2G or 3G network, and the UE needs to be handed over from a packet switched (PS) domain to a circuit switched (CS) domain, interruption may occur in a voice call, which affects continuity of the voice call of a user.\nCircuit switched fallback (CSFB) means fallback of UE covered by LTE to a network that has a CS domain when the UE processes a voice service. The US is triggered to be handed over from an LTE network to a CS domain of a 2G or 3G network, falls back, in a handover manner, to the network that has a CS domain, and processes the voice service in the network that has a CS domain.\nHowever, in the prior art, after the UE falls back from the LTE network to the 2G or 3G network by means of handover, a registered public land mobile network PLMN of the UE changes from a first PLMN of the LTE network to a second PLMN of the 2G or 3G network.\nWhen the UE returns to the LTE network from the 2G or 3G network, the UE may preferably camp on the second PLMN because the current registered PLMN of the UE is the second PLMN. As a result, the UE cannot return to the first PLMN in the LTE network, an LTE network operator of the first PLMN encounters a user turnover in a CSFB process, and operators cannot cooperate with each other."}
{"text": "1. Field of the Invention\nThe present invention is directed to a master cylinder, particularly for an anti-locking hydraulic brake system comprising means for determining the position of the master piston within the longitudinal bore of the master cylinder housing.\n2. Description of the Related Art\nA brake system including a master cylinder of the general type described above is disclosed in German Specification DE-OS 37 51 603. The master cylinder, through a brake conduit, is in communication with a wheel brake. The brake conduit can be blocked by means of an inlet valve. The wheel brake, through a by-pass provided by an outlet valve, is in communication with a reservoir. A pump delivers fluid from the reservoir into the working chamber of the master cylinder. During deceleration, the rotating pattern of the wheel is permanently monitored to immediately detect if a wheel tends to lock. Upon detection of a locking condition, the pump drive is actuated to cause the pump to deliver fluid into the master cylinder. By opening and closing the inlet and outlet valves, the pressure in the wheel brake can be so modulated that the brake torque attains a value corresponding to the friction between tire and road. The pressure fluid delivered by the pump into the working chamber of the master cylinder forces the master piston back into its initial position. To prevent the piston from being completely returned, it has previously been suggested that the position of the master piston be determined and the pump delivery be regulated in response to the determined position of the master piston. According to greater detail in the aforementioned German Specification, it is insured that the master piston, during a brake slip control, assumes a position in advance of its initial position.\nThe means for determining the position of the master piston, according to the German Specification, may comprise, for example, a plunger actuating a switch and being excited by a ramp on the master piston. According to another embodiment in which the master cylinder is actuated by a vacuum brake force booster, a switch is positioned in the housing wall of the brake force booster, and the switch is actuated by the booster wall.\nBoth of the above proposed solutions disadvantageously require special structures for the master cylinder or for the vacuum brake force booster. Moreover, there is a high possibility of malfunction because the proposed corrective means do not work in a contact-free manner with the consequence that substantial adjustment efforts are needed to insure precise positioning of the master piston during a brake slip control."}
{"text": "The present invention relates to novel coloring phthalide compounds and recording materials using any of the phthalide compounds as coloring component capable of yielding images which absorb the light in the near infrared region by thermosensitive recording, pressure-sensitive recording or by use of laser beams.\nA conventional pressure-sensitive recording material utilizes a chemical reaction between a colorless or light-colored leuco dye and a color developer capable of inducing color formation in the leuco dye to form colored images. More specifically, the pressure-sensitive recording material comprises (i) a coloring leuco dye sheet coated with a leuco dye which is dissolved in an organic solution and microcapsuled and (ii) a color developer sheet coated with a color developer for the leuco dye and a binder agent. The coloring leuco dye sheet is superimposed on the color developer sheet and pressure is applied thereto in such a manner that the microcapsules containing the leuco dye are ruptured so as to react with the color developer.\nA conventional thermosensitive recording material generally comprises a support such as paper, synthetic paper or plastic film, and a thermosensitive coloring layer comprising as the main components a leuco dye and a color developer cable of inducing color formation in the leuco dye upon application of heat thereto to form colored images. The application of heat for such image formation is carried out, for instance, by a thermal head, a thermal pen, laser beams or a stroboscope.\nThese pressure-sensitive and thermosensitive recording materials are widely used because recording can be carried out more speedily by using a relatively simple device, without the complicated steps such as development and fixing of images, as compared with the recording materials for electrophotography and electrostatic recording.\nConventionally employed representative leuco dyes for use in such recording materials are Crystal Violet Lactone and Leuco Crystal Violet which are for blue-coloring, and fluorane compounds having an anilino substituent at its 7 position which are for black coloring. Further, in the conventional thermosensitive recording material, colorless or light-colored leuco dyes having, for example, lactone, lactam, or spiropyran rings, are employed, and organic acids and phenolic acid materials are employed as the color developers for the leuco dyes.\nRecently optical character-reading apparatus have been developed, for example, for reading bar codes in bar code labels, and are used in a variety of fields. In these apparatus, light sources emitting light having a light wavelength of 700 nm or more, such as a light emit diode and a semi-conductor laser, are in general use. However, the above-mentioned leuco dyes, when colored, scarcely absorb the light in the near infrared region having a wavelength of 700 nm or more, so that a light source emitting light in the near infrared region cannot be effectively used when the above-mentioned leuco dyes are employed in the recording materials.\nIn order to improve this shortcoming of the conventional leuco dyes, leuco dyes capable of absorbing not only light in the visible region, but also the light in the near infrared region, have been proposed, for instance, in Japanese Laid-Open Patent Applications 51-121035, 51-121037 and 57-167979. The leuco dyes proposed in these Japanese Laid-Open Patent Applications, however, are insufficient in the power of absorbing the light in the rear infrared region for use in practice."}
{"text": "When an optical storage device such as a Digital Versatile Disk (DVD) drive or a Compact Disk Read Only Memory (CD-ROM) drive is instructed by the host computer to retrieve data from a disk, the device will in most cases need to move an optical pickup or head to a different radial position. A servo algorithm implemented in firmware is used to command radial movement to an actuator to which the head is attached. To complete the servo loop, the current radial position of the head must be known. Information is stored on an optical disk in concentric or spiral tracks and position information is found by counting track crossings with a track count block 100 as shown in the block diagram of a typical optical storage system 10 in FIG. 1 while various waveforms are illustrated in FIG. 2.\nThe conventional optical storage device 10 includes an optical pickup and associated circuitry 12 for receiving optical signals from a disk and generating corresponding RF signals 200, a read channel 14, a decoder 16 and an interface 18 for transmitting/receiving data and command signals to/from a host device (not shown). Additionally, the device 10 includes an envelope detector 20 coupled to receive the RF signals 200 and generate a signal T-B.sub.-- ENV 202 representing the difference between the amplitude of the top of the RF signal envelope 200 and the amplitude of the bottom of the RF signal envelope 200.\nIf the optical device 10 reads DVD disks, a differential phase detector 22 generates a DVD.sub.-- PES (position error signal) 210 from phase detection signals A, B, C and D received from the pickup 12. If the optical device 10 reads CD-ROM disks (either exclusively or in addition to reading DVD disks), a differential amplifier 24 and an analog-to-digital converter 26 generate a CD-ROM.sub.-- PES 206 from phase detection signals E and F received from the pickup 12.\nExemplary waveforms of the RF signal 200, the T-B.sub.-- ENV signal 202, CD-ROM.sub.-- PES 206 and the DVD.sub.-- PES 210 are illustrated in FIG. 2. Each of these signals is input into the track count block 100. The conventional optical storage device 10 typically employs a dual-phase track counting method and requires that the T-B.sub.-- ENV signal 202 be filtered by an RF filter 101 and \"sliced\" by a slicer 102 relative to a threshold 204, resulting in an RF.sub.-- SLICE signal 212. Additionally, the appropriate PES (depending upon the type of disk in the storage device 10) must be selected by a multiplexer (MUX) 108, filtered by a PES filter 110 and sliced by a slicer 112 relative to a threshold 205 to generate a PES.sub.-- SLICE signal 214.\nCount logic 104 causes a counter 106 to increment or decrement depending on the phase relationship of RF and PES slice signals 212 and 214. When the optical pickup 12 moves in one direction relative to the tracks, the PES.sub.-- SLICE signal 214 leads the RF.sub.-- SLICE signal 212 by 90 degrees; if the pickup 12 moves in the opposite direction, the PES.sub.-- SLICE signal 214 lags the RF.sub.-- SLICE 212 by 90 degrees. (Although only the phase of the PES signal 206 changes when the track crossing direction changes, because of physical imperfections in the disk, a sinusoidal component called runout may be superimposed on the seek profile. The effect will be that the observed track count appears as if it is opposite to the intended direction. For a track counting system to be accurate, changes in direction must be accurately accounted for; consequently, both the RF and PES signals must be used.)\nWhen a seek is to be performed, servo firmware loads the counter 106 with a tracks-to-go number representing the number of tracks between the current track and the target track. The firmware also asserts a seek-in-progress signal to the counter 106 allowing the counter 106 to count down in response to a signal from the count logic module 104 and output a track count signal 216 having a resolution of 1/4 track. Additional counter input signals are seek-direction and quadrature-direction which are used to change the track count direction for a given RF to PES phase relationship. The spiral-compensation signal, also an input to the counter 106, increments or decrements the track count each time an index occurs during a seek, determined by the disk's spiral direction.\nThe dual-phase track count method has worked well in CD-ROM applications but has numerous drawbacks when applied to DVD. Maximum seek velocities in DVD devices are much higher than typical CD-ROM seek velocities and can cause unreliable detection of the RF modulation. For example, at 500,000 tracks per second in a 1.times. continuous linear velocity (CLV) DVD device, the actual velocity of the media under the head can be as slow as 80% of nominal. Assuming that there are 26,175,000 bits per second and that one-fourth of the expected RF modulation period is needed to define a peak, then ##EQU1## channel bits are available to define the peak. Because the average wavelength is 4.7 channel bits, an average of only one peak will be available to be peak detected by the envelope detector. Taking resolution effects into account, the peak detected signal is likely to be smaller than when the modulation is slower. Therefore, at high velocities, many of the track crossings will be missed by the RF.sub.-- SLICE signal 212.\nAnother drawback is that DVD signals are noisy and incompatible with dual-phase. The depth of modulation in the DVD RF signal 200 is specified to be smaller than the corresponding signal in a CD-ROM device. The signal to noise ratio is thereby reduced at the input to the RF slicer 103 and the track count error rate is increased. The purpose of the filter 10 is to reject some of the noise. Because the signal of interest is a sine wave with a widely varying frequency, an optimum filter would be a tracking bandpass filter. However, conventional technology does not currently allow a practical implementation of such a filter; rather, a low pass filter is typically used but at high seek rates, the bandwidth is greater than necessary, thereby permitting low frequency noise to pass. While the bandwidth can be changed, intervention by firmware is required.\nAnother measure of the error rate margin is the relative spacing of the edges of the slice signals 212 and 214. The rising edge of the signal 214 should be in the center of the positive pulse of the signal 212 and every edge should be in the center of an opposing signal pulse. Thus, the duty cycle of the signal should be 50% for constant seek velocities. To meet this constraint, the slicer threshold should be halfway between the maximum excursions of the signals at the input to the slicers 102 and 112. Because the depth of modulation can vary from disk to disk, the RF slicer threshold must be calibrated for each. However, even when calibrated, the threshold will be non-ideal in certain regions of the disk and the slice signals 212 and 214 will have poor duty cycles. The problem is exacerbated by the sawtooth-wave form of the DVD.sub.-- PES 210.\nTo compensate, a firmware Kalman filter can be used to estimate the track count and correct errors. However, a firmware Kalman filter can generally only detect large, defect-caused errors and can fail to detect small, noise-induced track slips. The filter also is heavily dependent on the track counts being received at regular intervals. If the T-B.sub.-- ENV or PES thresholds are incorrect, the RF or PES slice duty cycle will not be 50% and the count steps will be irregular, necessitating a larger error detection threshold. It is especially difficult for the firmware filter to detect errors as the filter instruction rate is 88,000-176,000 instruction cycles per second while the maximum track crossing rate is greater than 500,000 tracks per second.\nMoreover, current systems use no information about any previous state of the system. Consequently, errors in the track count can resemble instantaneous changes in velocity even though such changes are physically impossible."}
{"text": "The race is on among today's employers to identify, acquire, and benefit from the many advantages offered by electronic computers and computer systems. Grocery stores use computers to track and manage inventories; banks use computers for near-instantaneous transactions with businesses, other lenders, and their customers; manufacturing plants use computers to provide automated control over assembly lines; and countless others are using computers to simplify, streamline, and otherwise improve their lives and/or businesses. More and more of these uses will certainly be identified as computer technology improves.\nIt has become clear that today's children, or tomorrow's workforce, will absolutely have to possess at least some computer skills if they are to be successful. This is a growing concern of parents and educators alike, and there is already talk of a “digital divide” that will unfairly separate children into two classes: 1) those who were fortunate enough to learn computer skills in their primary or secondary education, and who are prepared for the emerging workforce; and 2) those who were not so fortunate, and who will face an uphill battle in the workforce. The divide is usually not the fault of the children, the parents, or the educators, and as such, is unfair. Instead, the reason for the divide usually stems from the costs associated with providing computers for children to study.\nPrimary and secondary schools (e.g., K-12) typically face tight budget restrictions, and computers are not cheap. Furthermore, many schools cannot afford the additional costs involved in setting up a computer laboratory or network, such as routing additional wiring through the building, or installing power and environmental control equipment.\nCost, however, is not the only factor limiting the number of computers currently in our schools. Another problem is the limited availability of teaching materials that may be used with such computers. Without an extensive supply of lessons and teaching materials for the computers, it is difficult for an educator to justify the cost of acquiring them. There is a need for computer-based lesson plans for students, and teaching materials that allows teachers to effectively monitor, control, and guide the students through their computer lesson plans.\nThere are also additional problems faced by an educator or teacher who is considering implementing computers in the classroom. One such problem is theft and/or vandalism. For those schools that have undertaken the high cost of implementing a computer system or laboratory, there is a need for measures that can protect the school's equipment from being stolen or vandalized. Another problem is the space available in the classroom. The typical computer (desktop and laptop) is large enough to occupy most of the workspace (e.g., a student's desk) that is available to the student, making the individual computer less desirable. There is a need for an educational computer system that will allow one or more students to have individual computers to work with, but which will also minimize the space occupied by such computers.\nAnother problem is expandability. A school's enrollment typically varies from year to year, which makes expandability a serious concern for any school wishing to install a computer network or system. There is a need for an educator's computer network to be expandable. This desired expandability is not limited to hardware. An educator's lesson plan may also need to be both up-to-date with current events, and flexible to accommodate various teaching styles.\nAnother problem arises due to the fact that typical computers, such as laptop computers, must be deactivated in order to change power sources and/or exchange batteries. Such a deactivation would be disruptive to an educator's lesson plan, particularly if a number of students need to undergo such a process at different times throughout a lesson. Accordingly, there is a need for a teaching system that satisfies the needs identified above, but may also minimize the disruption to the lesson plans caused by changing power sources of computers.\nThese concerns, along with numerous others, may be alleviated using one or more embodiments of the present invention."}
{"text": "The invention relates to a method of deriving a dimension of a detail of an object from a data set of data values assigned to a multidimensional space and relating to the object. The invention also relates to a data processor for deriving a dimension of a detail of an object from a data set of data values assigned to a multidimensional space and relating to the object.\nA method and a data processor of this kind are known from the international patent application WO 97/13457.\nThe data set assigns data values to positions in the multidimensional space. The data set notably comprises density values which represent the spatial density distribution of the object. The known method enables the dimensions of a detail, particularly of a blood vessel, to be derived from the data set, notably density values of a patient to be examined. Determination of stenosis of a blood vessel, necessitates accurate determination of the width of the relevant blood vessel. To this end, the known method includes the acquisition of a density profile and a maximum density value of the detail is derived from the density profile. Subsequently, according to the known method there are determined edge points where the density values in the density profile amount to half the maximum density value. The width of the blood vessel is subsequently calculated from the distance between edge points. Even though the known method for the calculation of the width of the blood vessel takes into account the fact that blurring may be involved in the measured data set, it has been found that inaccuracies occur in the measurement of the width of the blood vessel nevertheless.\nInaccuracies in the determination of the width of blood vessels from data acquired by means of computed tomography are also discussed in the article xe2x80x9cEvaluating the potential and problems of three-dimensional computed tomography measurements of arterial stenosisxe2x80x9d by David S. Ebert et al in Journal of Digital Imaging 11 (1998), pp. 151-157.\nIt is an object of the invention to provide a method of deriving the dimension of a detail of an object from the data set with an accuracy which is superior to that of the known method. It is notably an object of the invention to avoid inaccuracies due to differences in the spatial resolution in different directions during the determination of the dimension.\nThis object is achieved by a method according to the invention which is characterized in that a preferential direction is chosen in the multidimensional space, the spatial resolution of the data set of data values in the preferential direction being higher than the spatial resolution of the data values in at least one direction other than the preferential direction, the dimension of the detail being derived from data values in the preferential direction.\nThe data set of data values essentially has the highest spatial resolution in the preferential direction. The resolution represents the smallest dimension of details in the object which can still be faithfully reproduced by the data set. The smaller the smallest faithfully reproduced detail, the higher the spatial resolution will be. Specifically, in comparison with other directions the least blurring has occurred in the preferential direction during the measurement of the data set. It has been found that notably inaccuracies are avoided which, if no further steps were taken, would occur when the spatial resolution of the data set differs in different directions in the object. For example, when the data set is acquired by means of an X-ray computed tomography method, the preferential direction is situated in the scanning plane. This is because it appears that in the case of X-ray computed tomography the highest spatial resolution by far occurs in the scanning plane. In X-ray computed tomography an X-ray source and an X-ray detector are rotated together about the patient in the scanning plane. During the rotation of the X-ray source and the X-ray detector, the patient on the one side and the X-ray source with the X-ray detector on the other side can be displaced relative to one another, notably in the longitudinal direction of the patient, so that the X-ray source and the X-ray detector travel along a helical path relative to the patient. In that case the scanning plane is shifted relative to the patient during the rotation; the scanning plane is notably shifted along the longitudinal axis of the patient. In various orientations of the X-ray source and the X-ray detector, relative to the patient to be examined, the patient is irradiated by means of X-rays and density profiles of the patient to be examined are acquired by measuring the X-ray absorption in different directions in the body of the patient to be examined. Density values in different positions in the body of the patient are reconstructed from the values of the X-ray absorption measured in different directions. When the data values in the preferential direction is utilized to derive the dimension of the detail, the dimension is particularly accurately derived from the data set. The effects of blurring during the determination of the dimension of the detail can thus be avoided to a significant extent. The method according to the invention is particularly suitable for the determination of the dimensions of details whose dimensions hardly differ in various directions. For example, the method according to the invention is particularly suitable for determining the dimensions of the cross-section of arteries of a patient to be examined, because arteries practically always have a substantially round cross-section. The invention notably offers the advantage that an accurate result is obtained for the dimension of the detail of the object, such as the width of the artery, when the resolution of the data set is high in one direction, being the preferential direction, and is low relative to the dimension of the detail in another direction.\nWhile the invention was described herein with an employment of an X-ray computed tomography imaging technique to acquire the data values, the invention may be employed with other techniques to acquire the data values, such as, for example, a magnetic resonance imaging technique or a 3D ultrasound imaging technique.\nThe data values are preferably acquired in one or more scanning planes. It has been found that the spatial resolution of the data values in such a scanning plane is substantially higher than that in directions outside the scanning plane. This means that the preferential direction is situated in the scanning plane or extends parallel to the scanning plane. The invention is particularly suitable for determining the dimensions of a cross-section of an elongate detail of the object. An artery in the body of the patient to be examined constitutes an example of such an elongate detail. The transverse plane of such an elongate detail is determined substantially perpendicularly to the longitudinal axis. According to the invention the dimension of a detail can be accurately determined by deriving the dimension from data values in the preferential direction. It has been found that the preferential direction is usually situated in the scanning plane in which the data values are acquired. This is the case notably when the data values are acquired by means of an X-ray computed tomography method. In case the transverse plane is situated parallel to the scanning plane, the dimension of the detail can be simply derived from data values in an arbitrary direction in the transverse plane, because the data values are blurred only slightly in essentially all directions in the transverse plane. The invention offers the advantage of accurate results concerning the dimension of the detail also in the case of data sets acquired by means of a magnetic resonance imaging system or a 3D ultrasound method.\nFurthermore, the dimension of the detail is preferably derived from data values relating to a perpendicular cross-section of the detail, thus avoiding, the involvement of an oblique projection, relative to the longitudinal axis of the detail, in the determination of the cross-section of the detail. To this end, the (local) longitudinal axis of the detail is derived from the data values in the data set; it is to be noted, however, that the longitudinal axis of the detail can also be derived from anatomical information. The transverse plane is situated substantially perpendicularly to the local longitudinal axis of the detail. The data values which are blurred the least and relate to the cross-section of the detail lie along the line of intersection of the transverse plane and the scanning plane. The width of the detail can thus be accurately determined by deriving the dimension from the data values along the line of intersection. The occurrence of inaccuracies in the determination of the dimension which are due to oblique cross-sections or blurring of the data values is thus avoided as well as possible.\nThe dimensions of the detail, such as the width of the artery, are accurately determined, for example by determining the locations where strong gradients occur in the data values in the preferential direction. Such gradients represent the edges of the artery; it is notably when the artery is filled with an X-ray-absorbing contrast medium during the acquisition of the data values by means of X-ray computed tomography that the density values in the artery are substantially higher than the data values of the surrounding tissue. Another approach for accurately determining the width of the artery is to determine the locations where the data values amount to a predetermined fraction, for example half, of the maximum density value in the detail, for example the artery. It has been found that the positions in which the density value amounts to half the maximum density value offer an accurate indication of the edges of the artery.\nThe method according to the invention yields an accurate dimension of the detail, for example the width of the artery. This result is a technical aid that can be used by a physician in making a diagnosis in respect of arterial stenosis in the patient to be examined.\nThese and other aspects of the invention are apparent from and will be elucidated, by way of example, with reference to the embodiment described hereinafter and the accompanying drawing; therein:"}
{"text": "Two classes of analysis systems are known in the field of medical analysis: wet analysis systems, and dry-chemical analysis systems. Wet analysis systems, which essentially operate using “wet reagents” (liquid reagents), perform an analysis via a number of required step such as, for example, providing a sample and a reagent into a reagent vessel, mixing the sample and reagent together in the reagent vessel, and measuring and analyzing the mixture for a measurement variable characteristic to provide a desired analytical result (analysis result). Such steps are often performed using technically complex, large, line-operated analysis instruments, which allow manifold movements of participating elements. This class of analysis system is typically used in large medical-analytic laboratories.\nOn the other hand, dry-chemical analysis systems operate using “dry reagents” which are typically integrated in a test element and implemented as a “test strip”, for example. When these dry-chemical analysis systems are used, the liquid sample dissolves the reagents in the test element, and the reaction of sample and dissolved reagent results in a change of a measurement variable, which can be measured on the test element itself. Above all, optically analyzable (in particular colorimetric) analysis systems are typical in this class, in which the measurement variable is a color change or other optically measurable variable. Electrochemical systems are also typical in this class, in which an electrical measurement variable characteristic for the analysis, in particular an electrical current upon application of a defined voltage, can be measured in a measuring zone of the test element using electrodes provided in the measuring zone.\nThe analysis instruments of the dry-chemical analysis systems are usually compact, and some of them are portable and battery-operated. The systems are used for decentralized analysis, for example, at resident physicians, on the wards of the hospitals, and in so-called “home monitoring” during the monitoring of medical-analytic parameters by the patient himself (in particular blood glucose analysis by diabetics or coagulation status by warfarin patients).\nIn wet analysis systems, the high-performance analysis instruments allow the performance of more complex multistep reaction sequences (“test protocols”). For example, immunochemical analyses often require a multistep reaction sequence, in which a “bound/free separation” (hereafter “b/f separation”), i.e., a separation of a bound phase and a free phase, is necessary. According to one test protocol, for example, the probe can first be transported through a porous solid matrix, which contains a specific binding reagent for the analyte. A marking reagent can subsequently be caused to flow through the porous matrix, to mark the bound analyte and allow its detection. To achieve precise analysis, a washing step must be performed, in which unbound marking reagent is completely removed. Numerous test protocols are known for determining manifold analytes, which differ in manifold ways, but which share the feature that they require complex handling having multiple reaction steps, in particular also a b/f separation possibly being necessary.\nTest strips and similar analysis elements normally do not allow controlled multistep reaction sequences. Test elements similar to test strips are known, which allow further functions, such as the separation of red blood cells from whole blood, in addition to supplying reagents in dried form. However, they normally do not allow precise control of the time sequence of individual reaction steps. Wet-chemical laboratory systems offer these capabilities, but are too large, too costly, and too complex to handle for many applications.\nTo close these gaps, analysis systems have been suggested which operate using test elements which are implemented in such a manner that at least one externally controlled (i.e., using an element outside the test element itself) liquid transport step occurs therein (“controllable test elements”). The external control can be based on the application of pressure differences (overpressure or low-pressure) or on the change of force actions (e.g., change of the action direction of gravity by attitude change of the test element or by acceleration forces). The external control is especially frequently performed by centrifugal forces, which act on a rotating test element as a function of the velocity of the rotation.\nAnalysis systems having controllable test elements are known and typically have a housing, which comprises a dimensionally-stable plastic material, and a sample analysis channel enclosed by the housing, which often comprises a sequence of multiple channel sections and chambers expanded in comparison to the channel sections lying between them. The structure of the sample analysis channel having its channel sections and chambers is defined by profiling of the plastic parts. This profiling is able to be generated by injection molding techniques or hot stamping. Microstructures, which are generated by lithography methods, are increasingly being used more recently, however.\nAnalysis systems having controllable test elements allow the miniaturization of tests which have only been able to be performed using large laboratory systems. In addition, they allow the parallelization of procedures by repeated application of identical structures for the parallel processing of similar analyses from one sample and/or identical analyses from different samples. It is a further advantage that the test elements can typically be produced using established production methods and that they can also be measured and analyzed using known analysis methods. Known methods and products can also be employed in the chemical and biochemical components of such test elements.\nIn spite of these advantages, there is a further need for improvement. In particular, analysis systems which operate using controllable test elements are still too large. The most compact dimensions possible are of great practical significance for many intended applications.\nU.S. Pat. No. 8,114,351 B2 discloses an analysis system for the analysis of a body fluid sample for an analyte. The analysis system provides a test element and an analysis instrument having a dosing station and a measurement station. The test element has a housing an (at least) one sample analysis channel enclosed by the housing. The test element is rotatable around an axis of rotation which extends through the test element.\nU.S. Pat. No. 8,470,588 B2 discloses a test element and a method for detecting an analyte. The test element is essentially disk shaped and flat, and can be rotated about a preferably central axis which is perpendicular to the plane of the disk shaped test element.\nKim, Tae-Hyeong, et al. “Flow-enhanced electrochemical immunosensors on centrifugal microfluidic platforms.” Lab on a Chip 13.18 (2013): 3747-3754, doi:10.1039/c3lc50374g, (hereafter “Kim et. al.”) discloses a fully integrated centrifugal microfluidic device with features for target antigen capture from biological samples, via a bead-based enzyme-linked immune-sorbent assay, and flow-enhanced electrochemical detection. This is integrated into a Centrifugal microfluidic discs, also known as “lab-on-a-disc” or microfluidic CDs.\nMartinez-Duarte, Rodrigo, et al. “The integration of 3D carbon-electrode dielectrophoresis on a CD-like centrifugal microfluidic platform.” Lab on a Chip 10.8 (2010): 1030-1043, doi:10.1039/B925456K, (hereafter “Martinez-Duarte et. al.”) discloses a dielectrophoresis (DEP)-assisted filter with a compact disk (CD)-based centrifugal platform. 3D carbon electrodes are fabricated using the C-MEMS technique and are used to implement a DEP-enabled active filter to trap particles of interest.\nEuropean patent application publication EP 2 302 396 A1 discloses an analyzing device includes: an operation cavity that is adjacent to a first reserving cavity retaining a sample liquid, in a circumferential direction of rotational driving; a connecting section provided on a side wall of the first reserving cavity to suck the sample liquid by a capillary force and transfer the sample liquid to the operation cavity; and second reserving cavities that are disposed outside the operation cavity in the circumferential direction of the rotational driving and communicate with the outermost position of the operation cavity through a connecting passage. The connecting section is circumferentially extended farther than the liquid level of the sample liquid retained in the first reserving cavity.\nUnited States patent application publication US 2009/0246082 discloses an analysis device comprising a separation chamber for separating a sample solution into a solution component and a solid component, a holding channel for holding a predetermined amount of the separated solid component, a mixing chamber connected to the holding channel, an overflow channel connected between the holding channel and the separation chamber, a sample overflow chamber into which the sample solution remaining in the separation chamber is discharged, and a joint channel connecting the separation chamber and the sample overflow chamber. After the separated solution component fills the overflow channel with priority by a capillary force, the separated solid component is transferred to the holding channel via the overflow channel, and a predetermined amount of the solid component is measured. The solid component in the holding channel is transferred to the mixing chamber by a centrifugal force, and simultaneously, the sample solution remaining in the separation chamber is discharged to the sample overflow chamber by the siphon effect of the joint channel."}
{"text": "One common style of prior art poultry house is the so called post-style poultry house. These post-style poultry houses are quite prevalent and typically were built using posts prior to the wide spread use of structural trusses. A typical post-style poultry house is between about 34 and 40 feet wide and about 300 to 500 feet long. The posts in such a house are usually spaced from one another and are used to support the roof of the poultry house.\nWhile a post-style house as just described is sturdy and long-lasting, in recent years changes in poultry harvesting technology have made this less than an ideal design for continued use. For example, traditionally once the poultry are ready to be harvested from the poultry house and to be taken to a processing plant, they are gathered up and placed in cages for transport. Currently, this often is being done manually by picking up poultry and putting them in the cages. More recently, automatic harvesting equipment is being developed which constitutes a self-propelled vehicle with equipment at the front end thereof for gathering the poultry and placing them in cages. Such machinery can be obstructed in its operation by the posts. As a result, there exists a strong economic incentive for replacing the post-style houses with a clear span-style poultry house.\nUnfortunately, to simply knock down existing post-style poultry houses and replace them with new, clear span poultry houses is prohibitively expensive for many poultry farmers. Moreover, it is economically wasteful inasmuch as much of the structure of the post-style poultry house might still be in good shape and should not be demolished and discarded. Accordingly, it can be seen that a need yet remains for a method and apparatus for economically converting post-style poultry houses to clear span poultry houses. It is to the provision of such a method and apparatus that the present invention is primarily directed."}
{"text": "This invention relates to a process for producing hydrocarbons useful as diesel fuel or a diesel fuel blending component from renewable feedstocks such as the triglycerides and free fatty acids found in materials such as plant oils, fish oils, animal fats, and greases. The process involves hydrogenation, decarboxylation and/or hydrodeoxygenation and isomerization in one or more reactors. A selective hot high pressure hydrogen stripper is used to remove at least the carbon oxides from the hydrogenation, decarboxylation and/or hydrodeoxygenation zone effluent before entering the isomerization zone. A diesel range stream or a naphtha range stream, or a mixture of the two streams is used as an additional rectification agent in the selective hot high pressure hydrogen stripper to decrease the amount of product carried in the overhead thereby reducing the amount of n-paraffins in the diesel fuel.\nAs the demand for diesel fuel increases worldwide there is increasing interest in sources other than petroleum crude oil for producing diesel fuel. One such non-petroleum source is what has been termed renewable sources. These renewable sources include, but are not limited to, plant oils such as corn, rapeseed, canola, soybean and algal oils, animal fats such as inedible tallow, fish oils and various waste streams such as yellow and brown greases and sewage sludge. The common feature of these sources is that they are composed of triglycerides and Free Fatty Acids (FFA). Both of these compounds contain n-paraffin chains having from about 8 to about 24 carbon atoms. The n-paraffin chains in the tri-glycerides or FFAs can also be mono-, di- or poly-unsaturated. Some of the glycerides from the renewable sources may be monoglycerides or diglycerides instead of or in addition to the trigylcerides. Another class of components considered a renewable feedstock is fatty acid alkyl esters.\nThere are reports in the art disclosing the production of hydrocarbons from oils. For example, U.S. Pat. No. 4,300,009 discloses the use of crystalline aluminosilicate zeolites to convert plant oils such as corn oil to hydrocarbons such as gasoline and chemicals such as para-xylene. U.S. Pat. No. 4,992,605 discloses the production of hydrocarbon products in the diesel boiling range by hydroprocessing vegetable oils such as canola or sunflower oil. Finally, US 2004/0230085 A1 discloses a process for treating a hydrocarbon component of biological origin by hydrodeoxygenation followed by isomerization.\nApplicants have developed a process which comprises one or more steps to hydrogenate, decarboxylate (and/or hydrodeoxygenate) and isomerize the feedstock. The performance of the isomerization catalyst is improved by removing at least carbon dioxide from the feed to the isomerization zone. The presence of oxygen containing molecules including water, carbon dioxide, and other carbon oxides may result in the deactivation of the isomerization catalyst. The oxygen containing molecules such as carbon dioxide, carbon monoxide and water are removed using a selective hot high pressure hydrogen stripper which additionally contains a rectification zone."}
{"text": "Fibrous structures, particularly sanitary tissue products comprising fibrous structures, are known to exhibit different values for particular properties. These differences may translate into one fibrous structure being softer or stronger or more absorbent or more flexible or less flexible or exhibit greater stretch or exhibit less stretch, for example, as compared to another fibrous structure.\nOne property of fibrous structures that is desirable to consumers is the Geometric Mean Elongation of the fibrous structure. It has been found that at least some consumers desire fibrous structures that exhibit a Geometric Mean Elongation of greater than 17.4% as measured according to the Elongation Test Method and/or greater than 15% as measured according to the Elongation Test Method if the fibrous structure exhibits certain absorbency properties and greater than 15.8% as measured according to the Elongation Test Method if the fibrous structure is a multi-ply fibrous structures. However, such fibrous structures are not known in the art. Accordingly, there exists a need for fibrous structures that exhibit a Geometric Mean Elongation of greater than 15.8% as measured according to the Elongation Test Method."}
{"text": "The present subject matter relates generally to a chest wound seal for first responders. More specifically, the present invention relates to a chest wound seal that stores compactly, includes: four-sided sealing; a one-way venting valve with no rigid components; and a one-way valve and vent channel resistant to blockage from clotting; and, additionally, a chest wound seal that may be thin; flexible; conformal to body contours; aggressively adhesive to contaminated surfaces (stick and stay); highly resistant to externally-originated particulate contamination; functional when covered with patient's clothing; and configurable as venting or non-venting.\nWhen an individual suffers a puncture wound to the chest, such as from being shot or stabbed, the wound often penetrates into the chest cavity so as to puncture the parietal pleura, visceral pleura, and sometimes even the lungs. Such a wound allows air to flow freely through the open chest wound into the chest cavity and pleural space each time the patient breathes.\nThe condition wherein air penetrates into the pleural space between the lung and the chest wall through a wound hole in the chest wall is an open pneumothorax. When an open pneumothorax occurs, the normal mechanism by which the lung expands is lost; i.e., the fluid adhesion of the pleural surface of the lung to the pleural surface of the chest wall. Thus, the affected lung does not expand normally when the patient inhales.\nA tension pneumothorax can occur when a one-way valve is formed by the wound that allows airflow into the pleural space while preventing airflow out. In a tension pneumothorax, each inhalation traps air in the chest, increasing pressure on the lungs and ultimately causing them to collapse. Additionally, the increasing pressure pushes important organs, such as the heart, major blood vessels, and airways, toward the center of the chest. This shifting can cause further compression of the lungs and may affect the flow of blood returning to the heart. These additional complications in a tension pneumothorax make it a life threatening condition that requires immediate treatment.\nFor patients with an open wound to the chest, the severity of the open pneumothorax that develops can be minimized by sealing the open wound via an occlusive bandage. For proper sealing, the bandage must be air- and liquid-tight, must be flexible enough to conform to all body contours, and must have an adhesive strong and persistent enough to keep the seal in place while subject to patient movement and various body fluids. The occlusive bandage may be used to eliminate air penetrating into the pleural space through the wound hole when the patient inhales.\nIn addition to developing a tension pneumothorax by air leakage through the chest wound, a patient with an open chest wound that is covered by an occlusive chest bandage may also develop a tension pneumothorax in a different manner. For instance, the bandage may have eliminated air from entering the pleural space through the wound hole, yet if lacerated, the lung will continue to leak air into the pleural space. Consequently, every time the patient inhales, more air becomes trapped in the pleural space, causing more pressure to be exerted on the lacerated lung. Again, this increased pressure on the lacerated lung forces the lung to collapse.\nTo successfully prevent a tension pneumothorax from developing, the air flow into the pleural space via the wound hole must be eliminated, but the air coming from the lacerated lung must be allowed to escape the pleural space via the wound hole.\nTherefore, a need exists for a compact, low profile, adhesive chest wound seal for use by first responders that includes a one-way valve that ensures effective venting of air and/or other fluids from the pleural cavity during exhalation while sealing a chest wound against air entry during inhalation, without risk of being dislodged during treatment of the patient.\nU.S. Pat. No. 5,478,333 to Asherman (“Asherman”) discloses a one-way valve apparatus installed in a chest wound seal to be used for treating a tension pneumothorax, as well as treat hemothorax (blood accumulation in the pleural cavity). The adhesive used on Asherman fails when exposed to bodily fluids such as blood or sweat. If the adhesive fails, then Asherman will not exhibit a 4-sided seal, rendering its air-ingress-during-inhalation-prevention function useless.\nIn addition, the design of Asherman will only allow air egress during exhalation if the valve apparatus is free of encumbrances, such as clothing, impinging on the valve exit end. Such interference will kink the valve shut, rendering it inoperative, and thereby turning Asherman into a simple occlusive chest wound seal. Further, Asherman will also not vent if the patient is lying upon the chest wound/vent, which will also kink the vent exit. Such a condition will prevent air ingress during inhalation, but will also prevent air/fluid egress during exhalation, a situation that can lead to formation of a life-threatening tension pneumothorax.\nU.S. Pat. No. 7,504,549 to Castellani et al. (“Hyfin”) discloses a one-way valve fabricated as an integral part of a thin, flexible chest wound seal that claims to prevent tension pneumothorax. Hyfin discloses that its drainage channel terminates at the peripheral edge of the chest wound seal. The lack of a full perimeter of adhesive (hydrogel) in combination with a thin and flexible backing sheet can allow the Hyfin to fold along its drainage channel and allow the two sides of that channel to contact each other during Hyfin application, causing the channel to close upon itself. The hydrogel is aggressively sticky, and once stuck to itself, cannot easily be separated to reopen the occlusion in the channel.\nU.S. Pat. No. 7,834,231 to Biddle et al. (“Biddle”) discloses a plurality of one-way valve apparati installed in a thin, flexible chest wound seal that claims to prevent tension pneumothorax. Biddle contains a one-way valve with a fibrous filter in line with the air flow path. The limited clear egress fluid path will be prone to blockage if the blood from the chest wound clots. Clotting will render Biddle a simple occlusive chest wound seal.\nIn addition, the one-way valve mechanism is subject to interference by externally originated particulate contamination, which can prevent valve closure during inhalation. If a valve fails to seal during inhalation, air can enter the pleural space through the chest wound and cause the formation of a life-threatening tension pneumothorax.\nU.S. Pat. Appl. No. US 2012/0046582 to Hopman et al. (“Hopman”)—discloses a one-way valve apparatus installed in a chest wound seal that claims to prevent tension pneumothorax. The Hopman valve is a rigid, bulky one-way valve apparatus. The Hopman valve is susceptible to externally originated particulate contamination, which can prevent valve closure during inhalation. If the valve fails to seal during inhalation, air can enter the pleural space through the chest wound and cause the formation of a life-threatening tension pneumothorax.\nThe one-way valve can be protected from external contaminants by replacing its removable cover, but with the cover in place, the vent is rendered inoperative, thereby turning Hopman into a simple occlusive chest wound seal. Such a condition will prevent air ingress during inhalation, but will also prevent air/fluid egress during exhalation, a situation that can lead to formation of a life-threatening tension pneumothorax.\nIn addition, the rigid, bulky one-way valve apparatus of Hopman protrudes approximately ⅝ of an inch above the top surface of the chest wound seal. This bulkiness ensures that fewer of the units can be carried in first responder kits than other thinner, more flexible chest wound seals.\nAccordingly, there is a need for a chest wound seal for first responders that stores compactly, includes: four-sided sealing; a one-way venting valve with no rigid components; and a one-way valve and vent channel resistant to blockage from clotting; and, additionally, is thin; flexible; conformal to body contours; aggressively adhesive to contaminated surfaces (stick and stay); highly resistant to externally-originated particulate contamination; functional when covered with patient's clothing; and configurable as venting or non-venting, as described herein."}
{"text": "The present invention relates to switch devices which each are turned on against a user's intention and electronic devices which each include such switch device.\nWristwatches are known which each illuminate its face or liquid crystal time display unit with a luminescence element such as a lamp or electronic luminescence element so as to allow for us to recognize the time even in a dark place.\nWristwatches of this kind include some which is provided with a lamp which is lighted up in response to the user turning on a specific switch, and some which are lighted up for a predetermined interval of time in response to the user's turning on a switch.\nAs shown in U.S. Pat. No. 5,612,931, wristwatches are known which include a posture detection switch which detects one motion of a user's gestures made when the user looks at the wristwatch to thereby light on a lamp.\nWhenever the user desires to illuminate the face of the wristwatch, however, the user must manipulate the switch. Thus, for example, when the user' hands are occupied with some baggage, the manipulation of the wristwatch is troublesome.\nFor instance, in a wristwatch which uses a posture detection switch which detects that the face of the wristwatch is inclined up at 12 o'clock for the user to look at the face to thereby light up the lamp, it is unnecessary to manipulate the switch to light up the lamp. However, even in that case, when the wristwatch takes such posture, the lamp would directly be lighted up, thereby inviting useless electric power consumption, and shortening the battery longevity.\nTherefore, it is an object of the present invention to provide a switch device whose operation is never performed without the user's intended motion, and an electronic device (wristwatch) which uses such switch device."}
{"text": "This section introduces aspects that may help facilitate a better understanding of the invention. Accordingly, the statements of this section are to be read in this light and are not to be understood as admissions about what is prior art or what is not prior art.\nConventional programmable devices, such as field-programmable gate arrays (FPGAs), are capable of communicating with external devices using different signaling schemes having various voltage levels, e.g., 1.2V, 1.5V, 1.8V, 2.5, and 3.3V. To handle incoming signals at these different voltage levels, conventional FPGAs are implemented with programmable input buffers that can be selectively programmed to operate at different voltage levels. In addition to a core power supply VCC and an input/output (I/O) power supply VCCIO, some FPGAs include an auxiliary power supply VCCAUX of 3.3V or 2.5V. For a programmable device having VCC, VCCIO, and VCCAUX power supplies, an input buffer might have three different buffer stages: a VCC-powered buffer stage designed to be programmably selected to handle low-voltage signaling (e.g., 1.2V), a VCCIO-powered buffer stage designed to be programmably selected to handle mid-voltage signaling (e.g., 1.5V and 1.8V), and a VCCAUX-powered buffer stage designed to be programmably selected to handle high-voltage signaling (e.g., 2.5V and 3.3V)."}
{"text": "The use of the internet as a location for discussions has influenced the desire to analyze such communications for useful information. Example internet communication mediums include discussion forums, instant messaging services, email, etc. In particular, businesses want to know what people are saying in these forums about their business(es) and/or product(s). One way for businesses to learn this information is to assign classifications to messages. For example, messages may be classified based on the sentiment expressed in the message. Summaries of the sentiment classifications can be reviewed to better understand the overall content of the messages."}
{"text": "The invention is concerned with a petroleum product and a process fop obtaining said petroleum product which may be used fop the formulation of a fuel fop an internal combustion engine, and the invention is also concerned with the product obtained by way of the process. The gas oils currently on the market which ape either in the form of motor fuels fop internal combustion engines with compression ignition (diesel type) of in the form of fuels ape most frequently products which have been refined and which contain about 0.3% sulphur (expressed by weight of sulphur). They ape usually obtained following a hydpodesulphuration treatment from a charge which can come from direct distillation of a crude petroleum of from a specific treatment of a crude petroleum (pyrolysis of distillation followed by pyrolysis of the fraction recovered during distillation, for example, of thermal cracking of catalytic cracking) usually containing at least 0.8% by weight of sulphur.\nIn some industrial countries, standards fop the content of sulphur have been set, or ape due to be set in the very neap future. These standards ape becoming increasingly strict, particularly in the case of gas oils which ape intended to be used as motor fuels. Thus, in France, particularly from 1995, the sulphur content of these gas oils will have to be at the most equal to 0.05% by weight (500 ppm) whilst the gas oils according to the standards in force have a sulphur content which can be as much as 0.3%.\nSimilarly, the gas oils used in France as fuels fop internal combustion engines must currently have a cetane index which is at least equal to 48, and gas oils used as fuels must have a cetane index which is at least equal to 40. Stricter standards ape to be expected in the neap future, particularly those concerning gas oils which ape used as engine fuels.\nMoreover, given the diversity of charges to be treated (crudes of various origin, and viscoreduction-, coking-, hydpoconversion-, distillation- or catalytic cracking crudes) to produce the gas oil it is desirable to be able to offer the refiner a process which is flexible and which is capable of adapting the products formed to demand and to satisfy future requirements as far as sulphur and nitrogen levels, cetane index, colour and content of aromatics are concerned.\nFinally, all existing processes such as hydrodesaromatisation or hydrocracking which enable petroleum products to be obtained which have a low sulphur content and which have a relatively high cetane index expend considerable amounts of hydrogen. By way of example, the process for hydrodesaromatisation of a direct distillation charge whose distillation intervals (ASTM D86) are 180.degree. C.&lt;T 5% &lt;300.degree. C., 260.degree. C.&lt;T 50%&lt;350, 350.degree. C.&lt;T 95%&lt;460.degree. C., expends 0.6 to 1.1% of hydrogen in relation to the charge, whilst the hydrocracking process requires more than 2% hydrogen in relation to its charge. Now, the hydrogen charge of the refinery which usually comes from the catalytic gas reforming unit tends to be inadequate in the light of the fact that gas oil standards are becoming increasingly strict, resulting in an increase in hydrotreatments.\nMoreover, these existing processes result in a petroleum product being obtained whose cetane index does not exceed 63, this latter only being attained at the cost of hydrogenation of the aromatic hydrocarbons of the charge, which reaction uses up hydrogen. (see table).\nIt is therefore desirable to offer the refiner a process which will enable him to produce a petroleum product which satisfies the various standards which will be applicable in the very near future, and after 1995, as far as sulphur content is concerned.\nThe technological plan is illustrated by the following patents: U.S. Pat. No. 4,985,139, EP-A-0 215 496, GB-A-1 006 949 and GB-A-943 239."}
{"text": "The present invention relates to a connector, and more particularly to a connector in which an upper cover is closed to an open rear opening of a plastic housing to press a plurality of wires against piercing sections of terminals located in passages provided on the plastic housing, so that insulating skins of the wires are cut to allow electrical connection of conductors in the wires with the terminals without the need of soldering, and time and cost for assembling the connector are reduced.\nWith the highly developed technologies, computers have been designed to execute more and more advanced functions. In other words, the central processing unit (CPU) inside each computer has incredibly high digital operating capability. Moreover, with the popularization of Internet, users demand for even faster upload and download speeds via the use of, for example, optical fibers, broad bands, asymmetric/asynchronous digital subscriber line (ADSL), etc. Under this condition, all other peripherals of the computer system, including various kinds of driving and driven elements, must also have a high operating speed corresponding to the high-speed CPU. For this purpose, it is very important for signal transmission cables, particularly connectors provided at two ends of the cables, to transmit signals between the CPU and the various driving and driven elements at a very high speed. It is also known that there are many devices, such as the universal serial bus (USB), the cable connectors, etc., being provided with an externally connected bus that supports a transmission rate as high as 12 Mbps with each port thereof capable of connecting up to 127 devices at the same time, as well as the functions of plug-and-play and hot-plugging.\nFIGS. 1 and 2 are perspective and sectioned side views, respectively, of a conventional connector to show the assembling thereof. As shown, the conventional connector includes a plastic core 1, a plurality of terminals 2, and a plurality of wires 3.\nThe plastic core 1 is internally provided with a plurality of equally spaced terminal holes that extend from a front to a rear end of the plastic core 1. A flat plate 11 is rearward extended from the rear end of the plastic core 1, and a plurality of terminal slots are formed on the flat plate 11 corresponding to the terminal holes. Each of the terminals 2 is in the form of a flat strip divided into a bent front conductive section 21 and a rear soldering section 22. The terminals 2 are provided on a main body thereof with barbs. Each of the wires 3 has a bare front end 31.\nTo assemble the terminal 2 to the plastic core 1, the conductive section 21 of the terminal 2 is inserted into one terminal hole from a rear end of the plastic core 1, so that the soldering section 22 of the terminal 2 is positioned in a corresponding terminal slot on the flat plate 11. When the terminal 2 is inserted into the plastic core 1, the barbs on the body of the terminal 2 are forced into the terminal hole to hook to inner wall surfaces of the terminal hole, making the terminal 2 fixedly located in the terminal hole. The soldering section 22 of the terminal 2 is exposed from the rear end of the plastic core 1 for the bare front end 31 of the wire 3 to overlap thereat. Thereafter, the bare front end 31 of the wire 3 is manually soldered to the soldering section 22 of the terminal 2 using soldering iron and soldering stick. Finally, a molding machine (not shown) is used to integrally mold a plastic case 4 around the wires 3 and the plastic core 1 to complete the connector.\nDuring the process of soldering the terminals 2 to the wires 3, since the terminals 2 have been previously inserted into the plastic core 1, the high temperature of the soldering iron tends to deform the plastic core 1 to therefore adversely affect the quality of the produced connector. Moreover, in the conventional connector, the terminals 2 are forced into and hooked to the plastic core 1 to thereby locate in the connector. To enable the terminals 2 to firmly hook to the plastic core 1, it is necessary to provide barbs on two sides of the terminals 2. In forming the terminals 2 with barbs, increased amount of scrap is produced when punching the material strips to form the terminals 2. In assembling the conventional connector, the terminals 2 must be accurately inserted into the terminal holes in the plastic core 1 with a properly controlled insertion force to avoid biased or deformed terminals 2. Moreover, an additional mold must be made to produce the plastic case 4 for integrally and tightly enclosing and fixing the wires 3, the terminals 2, and the plastic core 1 together. The mold would inevitably increase the manufacturing cost of the conventional connector to adversely reduce the competition ability of the connector in the international markets.\nIt is therefore tried by the inventor to develop a connector to eliminate drawbacks existed in the conventional connector, so that the connector may be quickly assembled to reduce the manufacturing cost thereof.\nA primary object of the present invention is to provide a connector in which an upper cover is closed to an open rear opening of a plastic housing to press a plurality of wires against piercing sections of terminals located in passages provided on the plastic housing, so that insulating skins of the wires are cut to allow electrical connection of conductors in the wires with the terminals without the need of soldering, and time and cost for assembling the connector are reduced.\nTo achieve the above and other objects, the connector of the present invention mainly includes a plurality of terminals, each of which includes a front conductive section and a rear piercing section; a plastic housing being provided at a front part with internal terminal slots for receiving the conductive sections of the terminals therein and at an open rear part with passages for receiving the piercing sections of the terminals; and an upper cover for closing the open rear part of the plastic housing. The upper cover is closed to the open rear opening of the plastic housing to press a plurality of wires against the piercing sections located in the passages on the plastic housing, so that insulating skins of the wires are cut to allow electrical connection of conductors in the wires with the terminals without the need of soldering."}
{"text": "1. Field of the Invention\nThe present invention generally relates to the field of helmet mountable displays (HMDs).\n2. Description of the Related Art\nU.S. Pat. No. 5,091,719 describes a helmet mountable display (HMD) in which left and right channels share a common optical path that includes first and second relay lenses having respective optical axes, and horizontal and vertical deflecting mirrors. The lenses are arranged so that their axes form a V, with the deflecting mirrors disposed at the apex. For the left (right) eye, an 8.times.1 array of fibers is modulated to project the intensity pattern of a 1k.times.1k image through the first (second) relay lens. A pair of motors oscillate the horizontal and vertical mirrors to deflect light through the second (first) lens, and raster scan the image onto a back projection screen which projects the image onto the helmet's visor.\nOne of the primary concerns in HMD design is the optical system's overall weight, which is dominated by the weight of the motors used to oscillate the mirrors. Although the described HMD uses a common optical path to reduce the number of lenses, it requires two motors to drive the respective mirrors. Another design consideration is that people who wear the helmet have different interpupillary distances (IPDs). The IPD of male pilots in the 5-95 percentile varies from about 55 mm to about 75 mm. In the `V` design, variations in IPD are compensated for by rotating the first and second lenses with respect to the apex point to achieve the correct spacing. However, this induces `parallelogram` distortion in the output image, i.e., a square input appears as a parallelogram, and requires electronic precompensation of the digital imagery.\nU.S. Pat. No. 5,166,778 discloses a single lens HMD that employs a single 8.times.1 array of fibers to modulate both left and right images through a single lens. Horizontal and vertical mirrors are oscillated by respective motors to reflect the intensity patterns back through the lens and onto respective back projection screens. The single lens design reduces the number of lenses, but the single lens must be very large to simultaneously provide images for the left and right eyes, and the system still requires two motors to scan the image. Variations in IPD are compensated by adjusting optical relays between the respective back projection screens and the visor.\nA single-eye head mounted projection display called \"The Private Eye\" is produced by Reflection Technology Corporation of Waltham, Mass. The system includes a column of LEDs which project successive columns of an image through a magnifying lens onto a mirror. The mirror is oscillated to horizontally sweep the projected patterns onto the pupil of the person wearing The Private Eye. The mirror must be relatively large and located near the eye to provide a reasonably large field of view. Additionally, for a two-eyed display, the system would require two drive motors. Since this system is one-eyed, IPD adjustment is not an issue."}
{"text": "Human speeches may be captured as signals representing speech or speech signals using transducers, e.g., microphones, and be further processed for a wide range of applications, e.g., noise reduction or denoising, speech recognition, speaker recognition, speech synthesis, hearing aids, cochlear implants, and speech compression. The speech signals may be captured in the form of analog signals or digital signals and may be stored in electronic media, e.g., memory or disks. The speech signals may be processed using analog or digital processors.\nOne commonly used scheme in speech signal processing may include transforming speech signals from a time domain representation, e.g., signal waveforms as a function of time, into a frequency domain representation, e.g., spectrums or spectrograms, using the Fourier transformation. However, the Fourier transform may have a fixed time-frequency resolution. Thus, frequency distribution is limited to a linear scale. This limitation may require additional processing, e.g., converting a linear frequency distribution to a non-linear frequency distribution, e.g., as in the basilar membrane. Additionally, the fixed window size of the Fourier transform may cause undesirable harmonics in spectrograms.\nThis problem may be illustrated by the following example. In most speech processing systems, speech signals are first converted into digital signals on which, for example, short-time Fourier transform via Fast Fourier Transforms (FFT) may be applied for computing spectrograms of the speech signals. The intensity of a spectrogram represents the amplitude of the signal at a particular time and at a particular frequency. FIG. 1 shows speech signals of a male voice recorded under two different scenarios. FIG. 1(A) shows the speech signal recorded using a close-talking microphone, and FIG. 1(B) shows the same speech signal simultaneously recorded using a hands-free microphone in a moving vehicle. The close-talking microphone is placed near the speaker's mouth, while the hands-free microphone is placed, e.g., on a sun visor of the vehicle. Due to the distance between the mouth and the microphone in the hands-free scenario, the recorded speech signal in FIG. 1(B) may include substantial background noise compared to the speech signal as shown in FIG. 1(A). Both speech signals in this example are recorded at a sampling rate of 8 KHz.\nFIG. 2 shows the spectrograms of speech signals shown in FIG. 1. Specifically, FIG. 2(A) shows a spectrogram of the clean speech signal recorded using a close-talking microphone as shown in FIG. 1(A), and FIG. 2(B) shows a spectrogram of the noisy speech signal recorded using a hand-free microphone as shown in FIG. 1(B). The spectrograms are computed in the standard way, as in traditional feature extraction for speaker or speech recognition. It may include using a window of, e.g., the length of 30 ms shifted every 10 ms with an overlap of 20 ms, of a type of Hamming window before applying FFT to the speech signals. To facilitate further analysis and feature extraction, the frequency distribution is mapped from linear scales to the Bark scale after the FFT computation.\nThe spectrograms of FIG. 2 shows two types of distortions: (1) the pitch harmonics as periodical waves along the frequency axis; and (2) the background noise in the frequency domain in the form of “snow” noise (which is different from environmental noise). Both types of distortions are caused by FFT computation.\nAdditional processing steps may be needed to remove effects of these undesirable distortions. FIG. 3 shows a schematic diagrams of such a Fourier transform based system 300 for extracting features from speech signals. A speech signal may first undergo pre-processing steps of frame blocking 302 and windowing 304. Then at 306 FFT may be applied to the pre-processed speech signal for computing a spectrum. As discussed above, the spectrum may include unwanted harmonics and noise. In this method, triangular filters 308, a non-linear filter 310, and discrete cosine transform (DCT) 312 may be applied to reduce the effects of harmonics and noises caused by FFT. However, while these steps may reduce the harmonic distortions, they may also remove useful and real pitch information. It is commonly known that humans identify peoples' voices by the voice characteristics, represented mainly by pitch information.\nSpeech signal processing based on wavelet transforms may provide flexible time-frequency resolution. However, it may have other noticeable problems. First, the existing wavelet does not have the capabilities of mimicking the impulse responses of the basilar membrane. Thereby, it may not be used directly for modeling the human hearing system, e.g., cochlea. Second, there is no discrete formula to approximate inverse continuous wavelet transforms even though forward and inverse transforms for continuous wavelet based are well defined."}
{"text": "A lithographic apparatus is a machine that applies a desired pattern onto a substrate, usually onto a target portion of the substrate. A lithographic apparatus can be used, for example, in the manufacture of integrated circuits (ICs). In that instance, a patterning device, which is alternatively referred to as a mask or a reticle, may be used to generate a circuit pattern to be formed on an individual layer of the IC. This pattern can be transferred onto a target portion (e.g. comprising part of, one, or several dies) on a substrate (e.g. a silicon wafer). Transfer of the pattern is typically via imaging onto a layer of radiation-sensitive material (resist) provided on the substrate. In general, a single substrate will contain a network of adjacent target portions that are successively patterned. Known lithographic apparatus include so-called steppers, in which each target portion is irradiated by exposing an entire pattern onto the target portion at one time, and so-called scanners, in which each target portion is irradiated by scanning the pattern through a radiation beam in a given direction (the “scanning”-direction) while synchronously scanning the substrate parallel or anti-parallel to this direction. It is also possible to transfer the pattern from the patterning device to the substrate by imprinting the pattern onto the substrate.\nIt has been proposed to immerse the substrate in the lithographic projection apparatus in a liquid having a relatively high refractive index, e.g. water, so as to fill a space between the final element of the projection system and the substrate. The point of this is to enable imaging of smaller features since the exposure radiation will have a shorter wavelength in the liquid. (The effect of the liquid may also be regarded as increasing the effective NA of the system and also increasing the depth of focus.) Other immersion liquids have been proposed, including water with solid particles (e.g. quartz) suspended therein.\nHowever, submersing the substrate or substrate and substrate table in a bath of liquid (see, for example, U.S. Pat. No. 4,509,852, hereby incorporated in its entirety by reference) means that there is a large body of liquid that must be accelerated during a scanning exposure. This requires additional or more powerful motors and turbulence in the liquid may lead to undesirable and unpredictable effects.\nOne of the solutions proposed is for a liquid supply system to provide liquid on only a localized area of the substrate and in between the final element of the projection system and the substrate (the substrate generally has a larger surface area than the final element of the projection system). One way which has been proposed to arrange for this is disclosed in PCT patent application publication WO 99/49504, hereby incorporated in its entirety by reference. As illustrated in FIGS. 2 and 3, liquid is supplied by at least one inlet IN onto the substrate, preferably along the direction of movement of the substrate relative to the final element, and is removed by at least one outlet OUT after having passed under the projection system. That is, as the substrate is scanned beneath the element in a −X direction, liquid is supplied at the +X side of the element and taken up at the −X side. FIG. 2 shows the arrangement schematically in which liquid is supplied via inlet IN and is taken up on the other side of the element by outlet OUT which is connected to a low pressure source. In the illustration of FIG. 2 the liquid is supplied along the direction of movement of the substrate relative to the final element, though this does not need to be the case. Various orientations and numbers of in- and out-lets positioned around the final element are possible, one example is illustrated in FIG. 3 in which four sets of an inlet with an outlet on either side are provided in a regular pattern around the final element."}
{"text": "Drilling fluid circulating systems are employed during downhole drilling operations to provide a number of functions. These functions include providing cooling to the drill bit, removing drilled cuttings to the surface and providing a hydrostatic head to the formation to prevent the escape of oil and gas from the well.\nAs is well known, in almost all drilling operations, drilling fluid is continuously lost to the formation during the drilling process due to seepage of the drilling fluid into the formations. Seepage will occur at varying rates depending on the relative porosity or fractures of the formation and the hydrostatic head pressure. While some seepage will always occur, generally it is necessary to control drilling fluid loss in order to ensure effective drilling fluid circulation and to reduce drilling fluid costs. Notwithstanding efforts to prevent drilling fluid losses, tens of thousands of cubic meters of these fluids are lost each year.\nThere are generally two classes of drilling fluids namely water-based and oil-based drilling fluids. Water-based drilling fluids are generally less expensive than oil-based drilling fluids but cannot be used effectively in all formations and may result in operational problems in certain formations such as hydratable shales, silts or clays. In these formations, there is a tendency for the hydratable materials to destabilize the wellbore as well as disperse within the drilling fluid. The dispersion effect causes substantial increases in the solids content of the drilling fluid and leads to various problems including solids separation difficulties at the surface and detrimental increases to the viscosity of the fluid.\nIn order to overcome these difficulties associated with water-based drilling fluids, drilling fluids comprised of oil or water-in-oil (a.k.a invert emulsions) are often used. Generally, emulsion compositions in which oil remains the predominant or continuous phase are preferred. In order to maintain effective and stable emulsions, the compositions will require the use of emulsifiers to stabilize the emulsion. Such emulsions are effective in stabilizing hydratable shale or clay formations.\nIt is also well known that seepage losses can be mitigated through the addition of fibrous materials into the drilling fluid such as sawdust, fabrics, nut hulls, seed hulls and minerals such as ground limestone. These materials can bridge the various loss zones where fluids are being lost by mixing with drilled cuttings, and “viscosifiers” to create a “filter-cake” that may significantly reduce fluid losses to the formation. Examples of the use of various plant derived components as drilling fluid additives include U.S. Pat. No. 6,323,158, U.S. Pat. No. 5,763,367, U.S. Pat. No. 5,599,776, U.S. Pat. No. 2,691,629, U.S. Pat. No. 6,750,179, U.S. Pat. No. 4,439,328, U.S. Pat. No. 5,076,944, U.S. Pat. No. 5,861,362, United States Patent 2004/0014609 and United States Patent Applications 2004/0063587 and 2004/00232816. In the past, any plant materials have generally been preprocessed to remove any oils and gums associated with the plant material.\nTypically, these materials also have a detrimental effect on the drilling fluid and may significantly change the properties of the fluid by increasing viscosity to unmanageable levels or de-stabilizing the emulsion. These effects directly impact the efficiency of surface pumping and solids removal equipment as well as causing wellbore destabilization. In the particular case of fibrous and granular materials including but not limited to cardboard, nut shells, calcium carbonate, and sawdust, within invert emulsions, these materials will be broken down over time to smaller sizes resulting in an increase in the relative surface area of these materials within the emulsion. The net effect of this size breakdown and the hydrophilic nature of these materials will contribute to the breakdown of the emulsion. As a result, over time, additional emulsifiers and oil-wetting agents must be added to the emulsions to maintain the emulsion properties, which again contribute to the cost of the drilling fluid. In the specific case of adding limestone or other mineral based seepage control agents to the drilling fluid, the specific gravity of the drilling fluid is also increased which increases the hydrostatic pressure within the wellbore which may directly increase the seepage losses and also requires stronger pumping equipment.\nOther functional requirements of drilling fluids, problems that may be encountered with the use of drilling fluids and past solutions are briefly discussed below:\nDrilling fluid systems also require the use of torque reducing and drag reducing agents to relieve both rotational twist that builds up in the drilling pipe during rotation and the frictional forces required to lift the drill string from the borehole through the drilling fluid. In the past, the use of spherical media such as glass beads, Teflon beads, styrene-divinylbenzene copolymer beads, walnut hulls, and oil based additives (petroleum or vegetative) within the drilling fluid have demonstrated the ability of these additives to relieve rotational twist and reduce drag.\nDrilling fluids also preferably require materials to disperse entrained drilled solids that build up in the drilling fluid system during the excavation process. In many systems, the materials used are surfactants and phenolic compounds including lignins and tannins.\nFiltration control agents may be added to both water and oil-based drilling fluids to reduce the penetration of the drilling fluid into formation rock media. The addition of colloidal, polymeric and colloid dispersing chemicals like surfactants, emulsifiers, lignite and/or lignosulphonate materials may be used.\nIn both oil and water based drilling fluids, emulsifiers may be added to aid the emulsification of oils added to water based fluids and water added to oil-based fluids. Emulsifiers also act as surfactants in both water based and oil based drilling fluids. Typically this effect is observed as a reduction in the viscosity of the treated fluid. The effect is caused by the adsorption of the emulsifier onto the surface of an entrained particle or emulsified phase component which enhances steric hindrance to flocculation between dispersed molecules. Generally, the addition of emulsifiers is dictated by the requirements of the fluid for the level of dispersion of secondary phase and/or entrained solids in the fluid system.\nIn the past, emulsifiers have been prepared using crude tall oil fatty acids (CTOFA) derived from Kraft pulp-making processes. CTOFA's are often used as the source of fatty acids for creating soaps/emulsifiers in oil-based drilling fluids. CTOFAs are normally considered to be a waste product and as a result, the price of CTOFAs is relatively low (in the range of $310/metric ton). The chemical composition of CTOFAs is primarily C18 oleic and linoleic acids in addition to relatively high concentrations of resins including abietic, dehydroabietic, palustric, isopalustric, isopimaric, neoabietic and pimaric acids. CTOFAs are very viscous and must typically be diluted with solvents in order to be efficiently mixed within a drilling fluid.\nSome invert emulsions are blended with purified phospholipids to assist in oil wetting and dispersion of hydrophilic materials (drilled cuttings, barite and the like) which become entrained in the fluid. Purified phospholipids are used sparingly due to the high cost (typically in the range of $1000/ton).\nIn comparison, oil seeds are characterized by a much wider range of fatty acids with chain lengths that typically range from C14-C24 and naturally contain phospholipids. Different oilseeds contain very different specific compositions of fatty acids and phospholipids. From a cost perspective whole or crushed oilseeds are also in the same range as CTOFAs. However, any purified oilseed fatty acid is expensive (in the range of $1000+/ton).\nWhile many solutions to the requirements of drilling fluids have been provided, there continues to be a need for enhanced drilling fluid compositions, while maintaining functionality to effectively bridge fractures and other loss zones. In particular, there has been a need for drilling fluid compositions having lipophilic/oleophilic characteristics to assist in dispersing the drilling fluid as the drilling fluid is broken down. Still further, there has been a need for drilling fluid compositions with increased stability where the stability of the drilling fluid is enhanced through delayed release of emulsifying agents."}
{"text": "The present invention relates to a process for spray-coating of mass-produced goods with frequent changes between at least one aqueous and at least one non-aqueous coating agent inside a spray-coating chamber.\nThe coating of vehicles in a production line takes place in fully automated vehicle coating lines whereby the application of the colour-determining coating layer (as finishing top coat layer or as base coat layer with a transparent top coal layer) is carried out with a number of different colour shades with frequent changes of colour. The change of colour or colour shade results from the manufacturing schedule and can take place after each coating operation of an automotive body, or the bodies are assembled in small groups of the same colour shade. Each coating agent with a different colour shade is supplied to the spray-coating devices from a coating supply reservoir via a dedicated circulation line connected thereto.\nTo obtain an automatic change of colour shades, it is sufficient to switch from one circulation line to another circulation line, so long as the feeding lines, the metering devices, as well as the spray-coating devices themselves, are automatically rinsed between applications of different colour shades. When changing colour in spray-coating chambers, in which non-aqueous coating agents of different colour shades are handled, this automatic rinsing is carried out with mixtures of water-immiscible organic solvents conventionally used in coating agents. However, the colour change in the spray-coating chambers, in which aqueous coating agents with different colour shades are handled, is followed by automatic rinsing with rinsing media which contain as main constituents water and water-miscible organic solvents conventionally used in aqueous coating agents, and which contain optional additional auxiliary agents in small quantities. The respective rinsing medium is also supplied from the supply reservoir to the colour changer via a dedicated circulation line and can be supplied from there to the spray-coating devices via the supply lines and the metering devices.\nThe time available for the actual change of colour, and therefore for the rinsing operation, is dependent upon the speed at which the vehicles are being conveyed through the spray chamber and on the distance between the automotive bodies. In practice, the conveyor speed and the distance are adjusted to the period of time required for the change of colour and the rinsing process, which totals for example 5 to 20 seconds. Quite apparently, productivity suffers with any delay, so there is a need to decrease the amount of time needed between application of different coating colours or colour shades.\nThere has been a focus in recent years, particularly in the vehicle manufacturing/assembly plants, to replace non-aqueous (i.e., solvent-based) coating systems with aqueous coating systems due to environmental considerations. The actual replacement of non-aqueous coatings with aqueous coatings, however, has been problematic, especially with basecoats and topcoats which provide, respectively, a particular colour or a desired property or effect to the article (e.g., an automobile) being coated. The need to apply both solvent based and aqueous based coatings adds considerable complexity to the coating process, particularly if the same spray chamber needs to be used to apply both types of coatings. In such situations, rinsing of the supply lines, metering devices and spray-coating devices themselves, can't be carried out with either the conventional rinsing agents used for solvent based coatings nor those conventionally used for aqueous based coatings alone. Instead, the two types of rinsing media need to be used in sequential rinsing steps and in the proper order. A change of colour from an aqueous to a non-aqueous coating agent, for example, would result in the following sequence:\ntermination of the application of an aqueous coating agent with the colour shade A; PA1 rinsing of the supply lines, the metering devices and of the spray devices with PA1 a rinsing medium which contains, as essential constituents, water and water-miscible organic solvents conventionally used in aqueous coating agents and which contains optional additional auxiliary agents in small quantities; PA1 rinsing with mixtures of water-immiscible organic solvents conventionally used in non-aqueous coating agents; and PA1 application of a non-aqueous coating agent with the colour shade B.\nThe same applies for a change of colour from a non-aqueous to an aqueous coating agent. This multiple step process, with an additional rinsing step, obviously results in lower productivity as it is impossible to realize such a sequence for each change of colour in the short time that is available for the rinsing operation. In addition to the decrease in efficiency and productivity, the consumption of rinsing media as well as the quantity and costs for waste disposal would increase. It is, of course, possible to utilize two separate spray apparatus in a single spray chamber, one dedicated to aqueous coatings and the other dedicated to non-aqueous coatings. However, this alternative clearly adds capital costs to the vehicle manufacturing process and is thus not a satisfactory solution.\nSolvent mixtures are known in the art which are said to be suitable for cleaning surfaces contaminated with solvent coatings or water-borne coatings. For example, U.S. Pat. No. 5,700,330 discloses a solvent mixture consisting of 15 to 35% by weight of one or more glycol ethers and of 65 to 85% by weight of one or more aliphatic ketones, which can be used for cleaning instruments that have been contaminated either by coating agents containing solvents or by water-dilutable coating agents. An improved solvent mixture is taught by DE A 198 23 852, which discloses a solvent mixture consisting of 20 to 65% by weight of one or more glycol ethers, 9 to 40% by weight of one or more aliphatic ketones, 10 to 35% by weight of one or more aliphatic monohydric alcohols and 5 to 20% by weight of one or more organic solvents, selected from the group of the N-substituted pyrrolidones and/or butyrolactone.\nIt is an object of the invention to provide a process for spray-coating mass-produced substrates, in particular automotive bodies and their parts, using aqueous as well as non-aqueous coating agents inside a spray-coating chamber. The process is intended to overcome the aforementioned problems and to be suited particularly for coating operations that require frequent changes in the coatings being applied. A solution to the problem is obtained by coating mass-produced substrates with aqueous and non-aqueous coating agents inside one and the same spraying chamber and by rinsing the installations to be rinsed of that particular coat-applying unit with a universal non-aqueous rinsing medium of a given composition each time the coating agent is changed. By \"installations to be rinsed\" are meant in particular the supply lines from the coat changer to the spray-coating devices, including the metering devices as well as the spray-coating devices themselves."}
{"text": "A vacuum cleaner is an electric cleaning apparatus which utilizes a motor to drive blades to rotate at a high speed, to generate a negative air pressure in a sealed housing, and thereby suctioning dusts into a dust bag. Vacuum cleaners can be generally divided into vertical type vacuum cleaners, horizontal type vacuum cleaners, and portable type vacuum cleaners and the like according to the structure. With vacuum cleaners become miniaturized and portable, the improvement in structure and volume of a vacuum cleaner motor, as a core component of a vacuum cleaner, also attracts lots of attentions.\nA conventional vacuum cleaner motor generally includes a stator iron core and a rotor, the rotor is fixed on an output rotating shaft, and a load blade is further fixedly arranged on the rotating shaft. Bearings of the rotating shaft are placed at an outer side of each of the load blade and the rotor, that is, at two ends of the rotating shaft, such bearing arrangement is not compact. To fix the bearings, one of the bearings is fixed by a motor hood, and the other bearing needs to be fixed by an additional end cover, thus the structure is complicated and the volume is large. Also, if this structure is employed, the distance between the bearings is generally set large, which may cause a significant vibration and a high noise when the motor operates, and adversely affect the normal use of the vacuum cleaner, and thus is apt to cause premature damage and aging of parts and components in the vacuum cleaner motor, and reduces the service life of the vacuum cleaner.\nFor avoiding the structural limitations of the conventional vacuum cleaner motor, a motor structure is disclosed in the conventional technology. In the motor structure, one bearing is employed, and the bearing is fixed to a middle part of the rotating shaft, and a load and a rotor are arranged at two ends of the rotating shaft. In this way, to fix the bearing, the bearing can be directly fixed to a motor cover, without additionally providing a component to fix the bearing, thus simplifying the structure of the motor, allowing the entire motor to be compact, and reducing its volume. However, for ensuring the stability and balance in operation of the motor, generally, the bearing is required to be made long, which may lead to a high material cost of the bearing, and an increased production cost of the vacuum cleaner motor.\nTherefore, in view of the above issues, it is necessary to provide a new rotor structure, to simplify the structure of the motor, reduce the overall volume of the motor, to enable the bearing to be close to a load as much as possible, reduce the wobble at shaft ends of the rotor, lower the noise, improve the service life of the motor, and also save the material, and reduce production costs."}
{"text": "The present invention relates to computerized color graphics, color reproduction, and electronic printing systems. In particular, the present invention relates to a method and apparatus for characterizing colorants such as inks and predicting the color appearance of placing one or more colorant layers on top of a substrate such as paper or film as occurs, for example, in printing a page and photography.\nComputerized color graphics systems and electronic printing systems are known in the art. Typically, they enable a user to produce a color image (any visual two-dimensional pattern including text, graphic line art, continuous tone images, etc.), and from that image to produce a picture which can be printed via a color reproduction system, for example, by producing color separation plates for offset printing. There has been much effort in the past to develop ways to accurately predict or simulate the appearance of such images when printed on a substrate, for example paper or film, using a number of colorants, such as inks, the prediction carried out without actually printing the images.\nPrinting may be carried out using halftoning, also called screening, which is the process of creating the illusion of a continuous tone (xe2x80x9cCT,xe2x80x9d xe2x80x9ccontonexe2x80x9d) image using an output (e.g., printing) device capable only of binary output (ink deposited or not deposited at any location on a substrate). For color printing, several images (xe2x80x9cseparationsxe2x80x9d) are produced in the primary colorants (typically inks) used to print in color, and printed over each other in a press. For typical four color printing, four images are produced in cyan (xe2x80x9cCxe2x80x9d), magenta (xe2x80x9cMxe2x80x9d), yellow (xe2x80x9cYxe2x80x9d) and black (xe2x80x9cKxe2x80x9d), and each of these images are halftoned. Usually, digital halftoning is used together with an imagesetter, laser printer, ink jet printer, digital film recorder, or other recorder output device.\nIt is desirable to be able to calculate accurately how a picture will look when one prints an image, including a halftoned continuous tone color image, with a certain technique on a certain substrate using a certain set of colorants (e.g., inks). A method of predicting the color appearance can be used for example to display a simulation of the color appearance on a computer display, or to print a simulation of the color appearance on a more easily accessible and cheaper printer as a proof of what is to be printed finally in production.\nIt is desirable to do this for both reflection printing on an opaque substrate, and transmission imaging on a transparent substrate.\nBy accurately calculating we mean a maximum deviation between predicted color and actual color of the order of 5 CIELAB Delta E units, and an average deviation of about 2 CIELAB Delta E units.\nVarious methods are known for calculating the color resulting from superimposing a set of colorant layers on a substrate. These methods can be divided in two groups. The first group are characterized by requiring printing and measuring a relatively large number of overprints of the colorants. That is, this group includes methods that for a particular set of colorants, a particular printing technique (e.g., offset printing on a particular imagesetter), a particular substrate type (e.g., paper, or film for a photographic transparency or print, textile, sheet of plastic, etc.), a particular substrate color (e.g., the color of the paper, or of the transparent film in the case of a transparency, or of the textile, or of the plastic sheet, etc.) and a particular order of printing the colorants, involve printing a relatively large number of overprints of the colorants in the set, for example as patches. These patches are measured with a spectrophotometer or calorimeter and the measurements are used to calculate any overprint of colorants in the set using mathematical techniques, for example, interpolation. If carried out well and carefully, these methods can lead to accurate results. Modem color management techniques such as COLORSYNC(trademark) (Apple Computer, Inc., Cupertino, Calif.) and the methods promoted by the International Color Consortium (ICC, see http://www.color.org) use such techniques. While these techniques can produce accurate results, and also work for halftone images, there are several drawbacks with such methods. One is that a large number of color patches of overprints need to be made. For example, the IT8.7.3 chart (American National Standards Institute [ANSI] Committee IT8 for Digital Data Exchange Standards) contains nearly a thousand patches for a four color output. Hence it is very difficult to characterize sets of more than four colorants, for example printing with six or seven colors (xe2x80x9cHI-FIxe2x80x9d printing including PANTONE(copyright) Hexachrome from Pantone, Inc., Carlstadt, N.J.). There also are applications where inks other than cyan, magenta, yellow and black need to be used. Another drawback is that a set of patches will only be useful for accurately calculating overprints using the particular colorants and the particular colorant printing order used in the patches. Changing one colorant in the colorant set or changing the order of overprinting typically requires redoing the whole job of printing the set of patches.\nAn additional problem occurs when using such characterizations with more than four inks in a typical modem color management workflow, such as the ICC workflow. The ICC standard uses look up tables and interpolation to predict the CIELAB values of a particular output device. The number of nodes in these lookup tables increases exponentially with the number of inks, and the amount of computer resources becomes too high to be feasible.\nAlso included in this first group of known methods are those that use the well known Neugebauer equations to predict the color of an overprint of colorants, and their derivative based on Yule-Nielsen and spectral Neugebauer equations (See H. E. J. Neugebauer: xe2x80x9cDie theoretischen Grundlagen des Mehrfarbenbuchdrucksxe2x80x9d, Zeitschrift fur wissenschaftliche Photographie, Photophysik und Photochemie, Band 36, Heft 4, April 1937; J. S. Arney, C. D. Arney:xe2x80x9d Modeling the Yule-Nielsen Halftone effect,xe2x80x9d Journal of imaging science and technology, vol. 40, No. 3, pp. 233-238, June 1996; R. Rolleston and R. Balasubramanian: xe2x80x9cAccuracy of Various Types of Neugebauer Modelxe2x80x9d, ISandT and SID\"\"s Color Imaging Conference: Transforms and Transportability of Color, pp. 32-37, 1993). These Neugebauer-like models still need the knowledge of the color of the overprints of the primary inks, so still requires measurements overprint combinations of the colorants and the gradation steps of the colorants. Also, Neugebauer equations-based methods are known not to produce accurate results.\nA second group of prior art methods are those that determine spectral characterizations of individual colorants that can be used for predicting the color of overprints of so-characterized colorants. These methods, for a fixed substrate and printing technique, involve making one or more printouts of each colorant on one or more substrates, measuring the prints, and out of this data extracting a set of one or more parameters for each colorant that can be used to calculate an overprint of each colorant. These methods thus have the advantage of not requiring producing a large set of overprints. One such prior art method uses the two-parameter Kubelka-Munk method which describes an ink wither by one spectral parameter, (K/S)(xcex), or by two spectral parameters, scattering S(xcex) and absorption xcex1(xcex), where xcex is the wavelength. See James H. Nobbs: xe2x80x9cKubelka-Munk Theory and the Prediction of Reflectancexe2x80x9d, Rev. Prog. Coloration, Vol. 15, pp. 66-75, 1985.\nDetermining the two colorant parameters involves measuring the spectrum on a bare substrate and on a black substrate, and solving the resulting equations. There also exist in the literature refinements on the two-parameter Kubelka-Munk theory that incorporate internal reflection, anisotropic scattering and other second order effects. However, these methods are applicable only for full (i.e., 100%) coverage of a layer of ink of a particular thickness, and/or are not applicable for variable ink coverages, for example, halftoning at less than 100% dot percentage. In addition, the colorant parameters they determine are neither substantially independent of the lowest substrate color nor are capable of incorporating dependencies on other ink layers. Thus these parameters are not easily applied to predicting the appearance of sequentially applied inks at less than 100% coverage wherein the substrate with the previously deposited inks may be regarded as a new substrate, with characterizable influence from layer to layer.\nIn order to so predict the appearance of overprints, it is advantageous to spectrally characterizing colorants by parameters that are substantially invariant for all substrates of the same type, and independent of the substrate color, in that, when there is no interaction between colorant layers, one can treat a substrate with one or more colorant overprints on it as a raw colored substrate of the same type but a different color, that different color dependent on the one or more colorant overprints.\nIn summary, prior art techniques that are capable of accurately calculating overprints of inks, including rasterized (i.e., halftoned) inks, require making a number of overprints of the colorants, so do include determining spectral parameters of a colorant that spectrally characterize the colorant. Furthermore, prior art techniques that do spectrally characterize the colorants are not applicable to less than 100% coverage, do not produce parameters that are substantially independent of substrate color, and do not produce parameters that can take into account interactions between the colorant layers.\nThus there is a need in the art for a method of characterizing colorants capable of accurately predicting the color of an overprint of the colorants without requiring measurement of overprints of the colorants and capable of predicting the color of overprints of halftones.\nThe Parent Invention describes such a method for determining and using colorant parameters that are substantially independent on the substrate color. The use of these parameters is described when there is no interaction between the layers. At least two parameters are needed for non-scattering parameters and at least three parameters are needed for scattering colorants. Such colorant parameters also may be used when there is some interaction between the layers of colorants. Parameters in such a case also may be determined that are capable of characterizing some interdependency of the colorants deposited, so that even when there is interaction between layers, one can treat a substrate with one or more colorant overprints on it as a raw colored substrate of the same type. Needed is a more detailed description of how the Parent Invention is applicable to this case of the colorant layers influenced each other. That is, the effect of any colorant being influenced by one or more of the inks printed before or after in an overprint.\nExamples of when the effect of a colorant (e.g., ink) may be so dependent on inks earlier printed on which the new ink is to be printed, include phenomena such as xe2x80x9cink trappingxe2x80x9d in which an ink is xe2x80x9ctaken upxe2x80x9d less by the substrate when another ink has first been printed onto this substrate. Such a phenomenon could occur, for example, in offset printing. Another case is when the parameters of one or more inks already deposited onto a substrate are influenced on a new ink being deposited on top of these previously deposited inks. A lower ink may xe2x80x9cspreadxe2x80x9d as a result of additional inks deposited on top, so that the area the lower ink covers is increased. This is called xe2x80x9cspreadingxe2x80x9d herein, and might occur, for example, with ink-jet printing. Other examples when one colorant may be dependent on one or more overlying or underlying colorant layers include an ink-re-transfer-like phenomenon which could occur, for example, in dye sublimation printing, and mixing of the inks, which can occur in many printing techniques.\nIn the cases of xe2x80x9ctrappingxe2x80x9d and xe2x80x9cspreadingxe2x80x9d, the perceived ink quantity (halftone coverage percentage) also may change from the amounts laid down because of such interdependencies.\nThus there is a need in the art for a method for characterizing the spectral properties of colorants overprinted at a coverage percentage which may be less than 100%, on a reflective or transmittive substrate, the characterizing being by a small number of colorant parameters that are a function of the coverage percentage and that are substantially invariant to the substrate color for all substrates of a particular substrate type and for a particular printing technique, such colorant parameters capable of substantially incorporate the dependency on the colorants deposited before and after. Also, there still is a need in the art for a method and apparatus for accurately predicting the color of printing a set of colorants at any coverage percentages on a reflective or transmittive substrate, the method not requiring making and measuring overprints of the colorants, and applicable to overprinting where each colorant is laid at less than 100% coverage, for example using halftoning.\nAn object of this invention is a method for characterizing the spectral properties of colorants such as inks when printed on a substrate of a particular type by a small number of parameters which are independent of the base substrate color and that can incorporate interaction between other colorants printed on the substrate. A further object of this invention is a method and apparatus which use such characterization for predicting the color of printing a set of colorants such as inks on a substrate. Other objects would be clear from the description below.\nThese and other objects of the invention are provided for in a method and apparatus for determining a set of colorant parameters that spectrally characterize a colorant in order to predict the color spectrum of an overprint of the colorant deposited using a printing technique at a coverage percentage on a substrate of a substrate type, the set of colorant parameters substantially independent of the base substrate color and capable of incorporating mutual dependencies of the colorant with other colorants deposited on the substrate. For the example of offset printing using halftoning, these colorant parameters are a function of the dot percentage and wavelength.\nAt least two parameters are determined for characterizing a non scattering colorant, and at least three are determined for a colorant that has scattering properties. When mutual interaction of colorant layers in an overprint is possible, at least six parameters are determined for characterizing a non scattering colorant, and at least eight are determined for a colorant that has scattering properties.\nIn one aspect of the invention, the method includes making a number of sets of prints of the colorant at a range of coverage percentages on different background colors, for example, different dot percentages for the case of offset printing. To incorporate mutual interaction, more parameters need to be determined, so additional set of prints are made. For example, three sets of prints are made on three backgrounds for non scattering colorants and on four backgrounds for scattering colorants and in addition, if necessary to obtain a sufficient number of sets of prints on different backgrounds to determine the parameters, one or more similar sets of printouts are overprinted with one or more overprinting colorants.\nThe spectra of the prints, for example, the reflection spectra in the case of offset printing, are measured, and equations in the set of parameters are formed and numerically solved for the parameters. The backgrounds are, in one embodiment, two lightly colored backgrounds, a medium colored background, and, for the case of four sets of prints, a dark colored background and any overprinting is done with a lightly colored colorant. In one implementation for offset printing, the four backgrounds are formed by pre-printing the backgrounds, for example, in the case of the lightly colored background using 25% halftone dots, for the medium dark background, using 50% halftone dots with black ink, and in the case of the dark background used for the case of a colorant that has a scattering component, using 100% black ink, the overprinting happens for instance with a 100% colorant in a second set of prints. For coverage percentages for which prints are not made, interpolation is used to determine the parameters for such colorant coverages.\nSome colorants, called derived colorants, can be defined by a recipe of concentrations of basic inks. In another aspect of the invention, the set of colorant parameters of a derived colorant is determined by making measurements on a set of dilutions of each of the basic colorants on several background colors on one or more substrates of the same substrate type, a set of dilutions in this sense possibly including the undiluted basic colorant. The prints that are made of each of the dilutions at different colorant coverages are the same as would be made to determine the colorant parameters of each of the dilutions of the basic colorants. Measurements of the resulting spectra are made. From these, coefficients, for example Kubelka-Munk coefficients are determined for the dilutions at the coverage percentages, and, using interpolation, Kubelka-Munk coefficients are determined for the concentrations specified in the recipe of the derived colorant. Then, using a formula for mixtures of colorants using the coefficients, for example, the Kubelka-Munk theory for mixtures modified for all the printed coverage percentages, the spectra that would result by printing the derived colorant on the different backgrounds at the different coverage percentages are calculated without actually having to print the derived colorant. From these calculated spectra, using steps analogous to the case of determining the spectral parameters directly, the spectral parameters of the derived colorant are determined.\nIn another aspect of the invention, a method for using the parameters for determining a single overprint of a colorant on a substrate also is disclosed. Using the result that the colorant parameters are substantially independent of substrate color and can incorporate the effects of interdependency with other colorants, a method is disclosed for determining the spectrum of several overprints of several colorants by sequentially applying the method for one colorant. One first determines the coverage percentages xe2x80x98belowxe2x80x99 and xe2x80x98abovexe2x80x99 each colorant, i.e., for a particular colorant, those on which the colorant is deposited and those deposited over the colorant, and from these percentages, one modifies the colorant parameters and then determines the color of the colorants when deposited sequentially on the background in much the same way as if the colorant parameters were substantially independent of each other (see the Parent Invention). the reflection spectrum of an overprint of the first colorant on the raw substrate. One then treats the resulting print as a new substrate of the same type but a different color (the color given by the first obtained reflection spectrum) and in the same way determines the reflection spectrum of overprinting a second colorant on top of the first. One continues this process until one has determined the reflection spectrum of all the overprints of the several colorants.\nAlso disclosed is an apparatus for using the colorant parameters to determine the spectrum of overprints on a substrate of a set of colorants at a set of coverage percentages. The apparatus comprises a first memory for storing the spectrum of the substrate, a logic unit with inputs specifying the coverage percentages and a set of outputs which are the values of the colorant parameters (as modified by interaction between the colorant layers) at the coverage percentages. When inter-layer interaction may occur, some dot gain values also are input for a range of coverage percentages. The outputs of the logic unit are coupled to a combiner unit, which also has an input coupled to the first memory. The combiner unit determines the color spectrum of the overprints on the substrate of the set of colorants at the set of coverage percentages. One embodiment of the apparatus described comprises the logic unit being implemented as one or more lookup tables and an arithmetic unit. Another embodiment described comprises the logic unit being implemented as one or more interpolators for determining by interpolation the colorant parameters at the required coverage percentages for each colorant from values of the parameters at a set of fixed values of coverage percentages and an arithmetic unit that incorporates the calculation of the average coverage percentage below and above each colorant and the modification of the colorant parameters by these coverage percentages. The combiner unit is implemented as one or more arithmetic units, each arithmetic unit having an associated background spectrum as an input and a set of colorant parameters at a coverage percentage as another input and determining the spectrum of the overprint of that colorant on a substrate having the associated background spectrum. A parallel implementation is described in which the number of sets of interpolators in the logic unit is the number of colorants and the number of arithmetic units in the combiner unit also is the number of colorants. A serial implementation also is disclosed which is applicable to the case of no interaction between th layers and in which the logic unit comprises a single set of interpolators for a single colorant and the combiner unit comprises a single arithmetic unit for determining a single overprint. The operation of such a serial implementation is that the logic unit and arithmetic unit determine one overprint at a time in sequence, starting with an overprint on the raw substrate, and continuing with the second overprint being on the raw overprint with the first colorant, etc., until all overprints have been determined. After a number of cycles equal to the number of colorant overprints, the output of the arithmetic unit is the spectrum of all the overprints.\nAlso disclosed is a method and apparatus for simulating the appearance of overprints of a set of colorants on a display using colorant parameters determined according to one or more embodiments of the method of the invention for such determination. One embodiment of the method includes the steps of determining the spectrum of the overprints according to one or more of the above embodiments for determining the spectrum of overprints, and from the spectrum, determining the CIE-XYZ values of the overprint using CIE illuminant-observer weightings, and converting these XYZ values to RGB values to drive a CRT (or LCD, etc.) monitor using matrix multiplication and, in some embodiments, also one-dimensional lookup tables. Another aspect of the method includes using a matrix multiplier, one dimensional lookup tables, a multidimensional lookup table and interpolation for converting the XYZ values to the device dependent color values needed to drive a proof printer. One embodiment of the apparatus includes an apparatus for determining the spectrum of the overprints according to one or more of the above embodiments, and multiplier adders for determining the CIE-XYZ values of the overprint using CIE illuminant-observer weightings, and a matrix multiplier and, in some embodiments, also one dimensional lookup tables to produce RGB values to drive a CRT monitor. Another aspect of the apparatus includes a matrix multiplier, one dimensional lookup tables, a multidimensional lookup table and a multidimensional interpolator for converting the XYZ values to the device dependent color values needed to drive a proof printer.\nThese and other objects, features, aspects and advantages of the present invention will become more apparent from the following detailed description of preferred embodiments with the accompanying drawings."}
{"text": "1. Field of Invention\nThe present invention relates to a warmer, and more particularly to a portable warmer, wherein a warmer bag is received in a casing to safely warm up the heat exchanging fluid in the warmer bag through a safety arrangement such that when the heat exchanging fluid is heated above an usable temperature, the safety arrangement automatically cuts off an electrical connection between the heating element and the external power source for preventing the heat exchanging fluid from being overheated.\n2. Description of Related Arts\nDuring winter time, most people would like to use some warmer apparatus, such as electric heaters, in order to keep themselves warm. However, such warmer apparatus requires an electric power and has a relatively bulky size that the user cannot carry the warmer apparatus everywhere.\nOne kind of small and portable heating device, called warm water bag comprises a sealing bag, which is made of waterproof material such as plastic, having an inner cavity and an opening. In order to use the water bag, the user has to fill up the hot water into the inner cavity of the bag. Therefore, the user is able to hold the water bag close to the body of the user. In other words, the water inside the cavity gradually dispenses its heat through the sealing bag, so that the user can put the water bag to part of body, such as hands or feet, to keep himself or herself warm.\nAn improved water bag is made that the water bag is adapted to be placed in a microwave oven to heat up the water inside the water bag. Likewise, a built-in heating element is mounted into the water bag to heat up the water inside the water bag.\nHowever, the traditional water bag is inconvenient for people to use. Since the hot water is gradually cooling down in a timely manner, the water bag cannot provide the warming ability after a short period of time. Then, the user must pour out the water inside the sealing bag and refill the hot water. In other words, the user must repeatedly refill the hot water frequently, which is very inefficient and inconvenient.\nIn the other hand, it is convenient for the user to heat up the water inside the sealing bag by microwave oven. However, it is relatively dangerous for the user when the water inside the sealing bag is overheated. Accordingly, few seconds of contacting with the hot water will cause a serious burn. The overheated water inside the sealing bag will also cause a serious explosion.\nLikewise, by electrically connecting the heating element with the external power source, the water inside the sealing bag can be gently heated up by the heating element. However, when the water is heated, the volume of the water is expanded. When the water expands at a predetermined volume larger than the capacity of the sealing bag, the sealing bag will be forced to be popped and the hot water will be leaked from the sealing bag. Since the heating element is electrically connected to the power source, the hot water leaking from the sealing bag may cause an electrical short circuit or even a serious electric shock to the user. In other words, the user needs to find more secure and safe way to keep warm by using the water bag."}
{"text": "The invention relates to a distilling and desalination apparatus comprising an evaporation chamber and a condensing chamber."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for forming film, a process for manufacturing a device, an electro-optical device, and electronic equipment.\n2. Description of Related Art\nManufacturing processes of electro-optical devices used in displays, display light sources, or the like, and electronic devices such as semiconductor devices, contain a step for placing material on a substrate to form a film on the substrate. Quality and functionality are intimately associated with the technique for placing material and the technique for forming a film. Therefore the techniques are important to improving the performance of each of the above-described devices.\nA technique for placing material on a substrate is a method of discharging liquid material as droplets via a nozzle provided in a discharge head (for example, Japanese Unexamined Patent Application, First Publication No. 11-274671). This droplet discharge method generates less waste in the consumption of liquid material, and more easily controls the amount and the location of liquid material placed on a substrate compared with techniques such as a spin coating method.\nA technique for forming a film on a substrate using a droplet discharge method requires accuracy in placing liquid material on a substrate, and also requires an improvement in the accuracy and stability of the form of a film after the liquid film has dried.\nIncidentally, regarding a technique for making liquid material into droplets to place it on a substrate, there is a technique whereby a bank is formed on a substrate, and droplets are placed in the region surrounded by the bank (Japanese Unexamined Patent Application, First Publication No. 11-271753, for example). In this technique, since the region where droplets are placed is defined by the bank, being a partition, the droplets (liquid material) are prevented from spreading while wet.\nIn recent years, progress has been made in regards to the high density of circuits comprising devices, so fine detail is required in the pattern. Furthermore, devices are required to be of low cost.\nIn a bank forming technique as mentioned above, a bank is formed by patterning using a printing method, or patterning by lithography. However, in the printing method, it is difficult to handle the fine detail of the pattern. Furthermore, in the lithography method, since part of the bank forming material is removed along with patterning, there is a lot of waste in the amount of material used.\nA technique for forming a film on a substrate using a droplet discharge method requires accuracy of placing liquid material on a substrate, and also requires fine detail in a film and an improvement in film properties after the liquid film has dried."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and a system for recovering an extra fluorescent material (phosphor) produced in a step for coating phosphor screens of color image receiving tubes, and purifying the fluorescent material for recycle and reuse thereof, and more particularly to a method and a system for easily and simply effecting the above described purification at the manufacturing location of such color image receiving tubes.\n2. Description of the Related Art\nColor image receiving tube (color picture tube, color television picture tube) is manufactured through steps of preparing fluorescent material (phosphor) slurries of red, green, and blue colors being three primary colors, respectively, and applying these slurries to predetermined dot-shaped or stripe-shaped positions on phosphor screen, thereby forming phosphor film. In this case, a large amount of extra phosphor slurries are produced, and when such extra phosphor slurries are dumped without any consideration, it is not only wasteful of resources and manufacturing cost, but also it is affected adversely in view of safeguard for environment. Typically, rare earth fluorescent materials used as red fluorescent material, for example, a fluorescent material represented by Y.sub.2 O.sub.2 S:Eu is expensive, so that reuse of the material after recovering the same is important for reducing the manufacturing cost of color image receiving tube.\nHowever, these phosphor slurries contain additives such as carbon (dag), pigments, polyvinyl alcohol, chromium compounds such as bichromates and dispersants in addition to fluorescent materials. Furthermore, there is such a case where inclusion of fluorescent materials having different colors deteriorate color-developing characteristics of the fluorescent material to be reused in case of recovery thereof. In order to remove these additives and the fluorescent materials having different colors, very complicated steps are required in recovery and reuse of fluorescent material. In this respect, any simple and efficient regenerating or recycling method by which a fluorescent material can be purified up to a satisfactory degree could have not yet been found.\nFor instance, Japanese Patent Application Laid-Open No. 51-131486 discloses a regenerating or recycling method for fluorescent materials (phosphors) wherein a recovered fluorescent material containing a photosensitive binder prepared by dispersing a bichromate into polyvinyl alcohol is treated with an oxidizing agent, and examples of the oxidizing agent include hydrogen peroxide, alkali peroxide, periodic acid and the salts thereof, perchloric acid and the salts thereof, hypochlorous acid and the salts thereof, peroxosulfuric acid and the salts thereof, bichromates and the like.\nIn Japanese Patent Application Laid-Open No. 51-149186, disclosed is a regenerating or recycling method for fluorescent materials wherein tannic acid is added to a recovered fluorescent material containing graphite and chromium compounds, the resulting mixture is heat-treated at a temperature of 400 to 600.degree. C., whereby the graphite is oxidized to vaporize away the same as carbon dioxide, and the remaining oxide is washed with water to remove the same.\nAmong the impurities contained in the above described recovered fluorescent materials, particularly, carbon exhibits black color so that it absorbs easily light rays, and therefore emission luminance of fluorescent material is remarkably reduced in the case where such carbon remains in the regenerated fluorescent material. Furthermore, since carbon is chemically stable, it is difficult to dissolve or decolor the recovered fluorescent material by reacting carbon with any material in accordance with any chemical treatment, and accordingly, it is required to separate carbon from the fluorescent material to remove the same in regenerating process.\nOther than a method for washing simply a recovered fluorescent material with water for separating and removing carbon from the fluorescent material, disclosed are a method wherein a recovered fluorescent material is treated with an oxidizing material such as periodic acid and the salts thereof, and sodium hypochlorite, and then the material so treated is washed with water (Japanese Patent Preliminary Publication Nos. 2-504160 and 2-504162), a method wherein a recovered fluorescent material is subjected to reduction treatment, and then the material so treated is washed with water (Japanese Patent Application Laid-Open No. 63-154784) and the like methods. Any of these methods is efficient for separating and removing carbon at a certain degree.\nIn recent years, however, with increase in strict demand for emission luminance of color image receiving tube, improvements in luminance of phosphor which decides the luminance of color image receiving tube are required, while regenerated and recycled fluorescent material has been used as a part of raw fluorescent material. For this reason, improvements in characteristic properties of the regenerated fluorescent material are increasingly required. For such recent requirement of improvements in characteristic properties of regenerated fluorescent material, the above described conventional methods are insufficient for the one for removing carbon from the regenerated fluorescent material.\nJapanese Patent Application Laid-Open No. 53-51190 discloses a method wherein for removing included sulfide fluorescent materials from recovered rare earth fluorescent materials, a specified amount of silver ion is reacted with the sulfide fluorescent materials.\nJapanese Patent Application Laid-Open No. 53-52052 discloses a method for regenerating fluorescent materials characterized in that recovered Y.sub.2 O.sub.2 S fluorescent materials are subjected to acid washing and water washing, and then the resulting slurry is treated with a solution containing polyvinyl alcohol and ammonium bichromate, or with a solution prepared by adding further acrylic emulsion thereto without drying the slurry. Moreover, as an improved method of the above described one, Japanese Patent Application Laid-Open No. 53-53588 discloses a method for treating a recovered fluorescent material with a solution prepared by adding further an anionic surfactant to the above described solution to prevent inclusion of other fluorescent materials.\nIn Japanese Patent Preliminary Publication No. 2-504160, a method wherein periodate is used for removing an organic polymer such as polyvinyl alcohol stuck to a recovered fluorescent material as well as a method wherein chromium compounds existing in the aforesaid fluorescent material are permitted to dissolve into water by the use of an aqueous solution of an oxidizing material such as sodium hypochlorite, ammonium peroxodisulfate, and hydrogen peroxide, whereby the resulting solution is removed are disclosed. Furthermore, Japanese Patent Preliminary Publication No. 2-504162 discloses a method for recovering organic contaminants from a recovered phosphor slurry without requiring a calcining or sintering step by employing the above described method.\nHowever, an individual method of the above described methods is insufficient for recovering an extra phosphor slurry produced in a step for coating a phosphor screen and removing a variety of substances contained in the slurry to regenerate the phosphor, while a mere combination of these methods only brings about complicated steps. Particularly, organic materials such as polyvinyl alcohol and carbon cannot be completely removed, so that a fluorescent material having reusable characteristics cannot efficiently be obtained. For this reason, luminance and grade of coating film in a phosphor screen manufactured by employing the regenerated fluorescent material are adversely affected.\nFurthermore, there arises a necessity for regenerating or recycling a recovered phosphor slurry in a image receiving tube manufacturing plant without spending labor and cost for long distance transportation of such slurry instead of conducting concentratively regenerating treatment of the recovered phosphor slurry. In this respect, a simple regenerating method of fluorescent material for satisfying the above described necessity is demanded, but the purpose therefor could have not been attained by the combination of conventional methods as mentioned above.\nOn the other hand, the regenerating treatment as described above is carried out usually by such a manner that a recovered phosphor slurry is successively agitated, admixed, and washed with each of aqueous solutions of the above described chemicals in a variety of agitation tanks. Although it is also possible to treat such recovered phosphor slurry at normal temperatures, the treatment is mostly conducted with the use of hot aqueous solutions for elevating washing effect. For this reason, a mild steel material on the surface of which had been coated with polyvinyl chloride, polypropylene, or had been applied glass lining has been heretofore used for a tank body for the agitation tank and a baffle (baffle plate) disposed in the tank for elevating agitating efficiency, and such baffles have been fixed to the tank body by means of welding or bolting.\nSuch conventional agitation tank as described above has involved the following problems.\n(1) Baffles have been prepared separately from tank bodies and the former has been mounted on the latter, and in this connection, the working and mounting therefor are troublesome.\n(2) Mounted portions of baffles are liable to be loaded, and such portions are easily damaged.\n(3) Fluorescent material is easily stuck to and deposited on the mounted portions of baffles.\n(4) Protecting tubes for protecting a drain tube, a level sensor and the like to be inserted into a tank from agitation flow are necessary, and mounting for such protecting tubes are troublesome as same as in the case of mounting for baffles.\n(5) Agitating effect for liquid into which has been dispersed a fluorescent material is not sufficient.\nIn addition to the above described problems, polyvinyl chloride used for a surface coating of mild steel has a poor heat resistance, so that the polyvinyl chloride surface coating is corroded with hot aqueous solution. On one hand, polypropylene exhibits poor workability. Further, glass lining is easily attacked by aqueous alkaline solution in addition to exhibiting poor impact resistance. Other problems than those described above are such that sections to which such surface coating as described above is applied hardly are corroded with aqueous chemical solutions to be used, and the like problems."}
{"text": "The human proenkephalin A gene (PENK-A) contains 4 exons and codes for a series of structurally related oligopeptides like methionine enkephalin (Met-ENK), leucine enkephalin (Leu-ENK), methionine enkephalin arginine phenylalanine (Met-ENK-Arg-Phe) as well as methionine enkephalin arginine glycin leucine (Met-ENK-Arg-Gly-Leu), wherein one molecule of proenkephalin comprises the sequences of 4 Met-ENK and one Leu-ENK, Met-ENK-Arg-Phe and Met-ENK-Arg-Gly-Leu each. Met-ENK-Arg-Phe and Met-ENK-Arg-Gly-Leu are further metabolised to Met-ENK. The enkephalins have a function as neurotransmitters as well as neuromodulators and neurohormones. PENK-A besides the enkephalins comprises the sequences of enkelytin and peptide B at its N-terminus as well as synenkephalin at its C-terminus, that have an antibacterial function.\nThe expression of proenkephalin occurs in the central nervous system as well as the peripheral nervous system. In the brain increased concentrations of PENK are observed in the nucleus caudatus and the nucleus accubens, in the peri-aqueductal grey as well as in the hippocampus and the raphe nuclei.\nEnkephalins play a major role in a variety of physiological processes like perception of pain, regulation of stress response comprising hormone regulation, regulation of bone formation as well as regulation of immune responses. Met-ENK stimulates the proliferation of B- and T-lymphocytes and Leu-ENK the proliferation of T-helper and cytotoxic T-cells. Due to the immune regulatory properties of Met-ENK this enkephalin is classified as cytokine. Endogenous enkephalin peptides are also involved in normal regulation of cardiovascular functions like heart frequency, contractile strength and arterial blood pressure.\nEnkelytin has antibacterial function against gram-positive bacteria like Streptococcus aureus, Bacillus megaterium and Micrococcus luteus, however it is not inhibiting the growth of gram-negative bacteria like Escherichia coli, as well as the growth of fungi. Synenkephalin shows antibacterial properties against gram-positive bacteria as well as gram-negative bacteria.\nThe level of enkephalin is altered in body fluids and tissue in a variety of diseases. The Met-ENK immune reactivity in plasma of patients suffering from a brain infarctus in the acute phase and diabetes patients with symptomatic myocardiac ischemia is significantly increased as compared to healthy control individuals, as well as in plasma and liquor of patients suffering from migraine and headache and in liquor of schizophrenia patients.\nIn Parkinson's disease the content of Met-ENK-Arg-Gly-Leu in liquor of patients is significantly decreased as compared to healthy control individuals.\nPatients suffering from Alzheimer's disease showed a 4-fold increase in Met-ENK-like immune reactivity.\nThe concentration of enkephalin m-RNA in the striatum of patients suffering from Huntington's disease is significantly reduced, in liquor a decrease of met-ENK-Arg-Gly-Leu concentration was shown.\nIn tumor diseases prolactin-(PL) and adrenocorticotropin (ACTH)-secreting adenomas of the hypophysis show an increased concentration of Met-ENK as well as a 10-fold increased concentration of proenkephalin peptides in the latter.\nEnkephalins play a role in the pathophysiological response of nerve tissue in traumatic injuries like head injuries and injuries of the skin for example during and after surgery.\nEnkephalins furthermore play a major role in the pathogenesis of systemic inflammations. In animal models the induction of a peritonitis or the application of lipoypolysaccarides lead to an increase in PENK-expression as well as to an increase in Met-ENK-plasma concentrations.\nIn abscesses of periarthritic cattle different PENK-fragments could be detected, that comprise the sequence of enkelytin. In human peripheric monocytes and rat lymphocytes the induction of PENK-m-RNA expression by lipopolysaccharides could be demonstrated:\nThe infection of rats by borna-viruses resulted in an increased PENK-A transcription in the striatum.\nMet-ENK, Met-ENK-Arg-Phe and the proenkaphaline peptides B and F are significantly increased in plasma during hemorrhagic hypotension.\nThe biosynthesis of enkephalins (see FIG. 2) occurs as is the case in other peptide hormones, as preprohormone at (Golgi-) bound ribosomes. After separation of the hydrophobic N-terminal signal sequence by so-called signal peptidases and folding of the proteins in the lumen of the endoplasmic reticulum, the propeptides are packed into vesicles in the Golgi apparatus and are transported to the cell membrane. During transport the propeptides are processed to mature hormones by so-called prohormone convertases at usually dibasic amino acid sequences. Via different stimuli the peptides are secreted into the extracellular space or into the plasma. The mature peptides are rapidly inactivated after secretion by proteolysis.\nEnkephalins have a half-life of 12 to 15 minutes in plasma (ex vivo).\nThe processing of proenkephalin occurs in several steps in a strict timely order and is tissue-specific. An earlier step in processing is the separation of the C-terminal peptide B that comprises the sequences of enkelytin and Met-ENK-Arg-Phe. The processing of proenkephalin leads to the formation of different proenkephalin peptides.\nThe presence of the PENK-fragment 119-159 in liquor, however not in blood, was already shown by Stark et al., 2000.\nAs this fragment plays a central role in this invention, it is called “enkephalin-footprint” (ENK-fp) herein."}
{"text": "The sales of industrial robots that has been driven by the automotive industry is now moving into tasks as diverse as cleaning sewers, detecting bombs, and performing intricate surgery. The number of units sold increased to 120,000 units in 2010, twice the number as the previous year, with automotive, metal, and electronics industries driving the growth.\nPrior approaches to calibrating robots use measuring devices that either measure the inaccuracies of the robots after the robot is built or devices which measure work piece positions relative to the robots position prior to OLP's.\nSome of the prior art includes:    U.S. Pat. No. 7,979,159 (Fixell) discloses an invention which relates to a method and a system for determining the relation between a local coordinate system located in the working range of an industrial robot and a robot coordinate system. The method includes attaching a first calibration object in a fixed relation to the robot and determining the position of the first calibration object in relation to the robot. Then, locating at least three second calibration objects in the working range of the robot, a reference position for each of the second calibration objects in the local coordinate system can be determined by moving the robot until the first calibration object is in mechanical contact with each second calibration object. By reading the position of the robot when the calibration objects are in mechanical contact the relation between the local coordinate system and the robot coordinate system can be calculated.    U.S. Pat. No. 7,945,349 (Svensson, et. al.) discloses an invention which relates to a method and a system for facilitating calibration of a robot cell. One or more objects and an industrial robot perform work in connection to the objects, wherein the robot cell is programmed by means of an off-line programming tool including a graphical component for generating 2D or 3D graphics based on graphical models of the objects. The system comprises a computer unit located at the off-line programming site and configured to store a sequence of calibration points for each of the objects, and to generate a sequence of images including graphical representations of the objects to be calibrated and the calibration points in relation to the objects, and to transfer the images to the robot, and that the robot is configured to display said sequence of images to a robot operator during calibration of the robot cell so that for each calibration point a view including the present calibration point and the object to be calibrated is displayed to the robot operator.    U.S. Pat. No. 7,756,608 (Brogardh) discloses a method for calibration of an industrial robot including a plurality of movable links and a plurality of actuators effecting movement of the links and thereby the robot. The method includes mounting a measuring tip on or in the vicinity of the robot, moving the robot such that the measuring tip is in contact with a plurality of measuring points on the surface of at least one geometrical structure on or in the vicinity of the robot, reading and storing the positions of the actuators for each measuring point, and estimating a plurality of kinematic parameters for the robot based on a geometrical model of the geometrical structure, a kinematic model of the robot, and the stored positions of the actuators for the measuring points.\nWhat is needed are automatic and manual robotic tool finder systems and methods to improve cost and time factors in applications where absolutely accurate robots are not really necessary; examples including body-in-white (BIW) applications, resistance welding, material handling, and MIG welding.\nWhat is needed are automatic and manual robotic tool finder systems and methods for using different robot tools on a shop floor without having to perform a recalibration for each tool.\nThe primary objective of the automatic and manual robotic tool finder systems and methods of the present invention is to provide a calibration method which is simpler to operate, results in improved precision, involves a lower investment cost, and entails lower operating costs.\nAnother objective of the automatic and manual robotic tool finder systems and methods of the present invention is to increase the accuracy of the off-line program and decrease robot teaching time, while also negating the need for the technician to “jog” the robot into position as the process is automated."}
{"text": "The present invention generally relates to an oil pump. More specifically, the invention is an oil pump with dual scavenging for a twin cam engine.\nIt is an object of the invention to provide an oil pump with dual scavenging for a twin cam engine, with more than one pump returning oil from a sump or oil gathering location to an oil storage area or an area away from the engine.\nIt is an object of the invention to provide an oil pump with dual scavenging for a twin cam engine, with a design that improves oil control within the engine, which improves engine performance by maintaining a dry sump.\nIt is an object of the invention to provide an oil pump with dual scavenging for a twin cam engine, with a separate oil passageway from the crankcase to a first suction pump and another separate passageway from the cam chest to a second suction pump.\nWhat is really needed is an oil pump with dual scavenging for a twin cam engine, with more than one pump returning oil from a sump or oil gathering location to the oil storage area or an area away from the engine, with a design that improves oil control within the engine, which improves engine performance by maintaining a dry sump and a separate oil passageway from the crankcase to a first suction pump and another separate passageway from the cam chest to a second suction pump."}
{"text": "1. Field of the Invention\nThis invention relates to integrated circuit structures and processes for manufacturing them, and in particular to an improved power metal oxide semiconductor field effect transistor (MOSFET) structure having higher switching speeds and improved breakdown characteristics.\n2. Description of the Prior Art\nContinuing improvements in integrated circuit fabrication technology have lead to the manufacture of MOSFET's which are capable of efficient high frequency switching of relatively large amounts of power. The power MOSFET represents one of the first power semiconductor devices manufacturable using integrated circuit processing technology. A survey article discussing power MOSFET's is entitled \"Switching Lots of Watts at High Speed\" by B. Baliga, IEEE Spectrum, December 1981, pages 42-48.\nOne promising technology for the manufacture of advanced power MOSFET structures is VDMOS, meaning MOS devices fabricated using double diffusions in V-shaped grooves. Unfortunately conventionally fabricated VDMOS structures suffer from several disadvantages. Such devices are generally overly susceptible to secondary breakdown of parasitic bipolar transistors inherent in the structure. Furthermore, such structures typically have undesirably large drain to source resistance when on. In addition, the devices switch at undesirably low speeds in view of their relatively large size and therefore high gate capacitances."}
{"text": "1. Field of the Invention\nThe present invention relates to a structure of an engine. More particularly, the present invention relates to a structure of a gasoline direct injection engine having improved operational performance.\n2. Description of Related Art\nIn general, a GDI (gasoline direct injection) technology has been developed to improve fuel efficiency and performance of an internal combustion engine. The GDI engine technology is to directly inject fuel into the combustion chamber, not into the intake manifold.\nSince it is possible to directly inject fuel into the combustion chamber and produce a fuel-air mixture layer by using the GDI engine, it is possible to produce a condensed mixture by concentrate air and fuel around a spark plug. Accordingly, the engine can operate at a very small air/fuel ratio and wall wetting is reduced in comparison to injecting fuel to the intake port in the related art, such that it is possible to accurately control the amount of fuel and improve fuel efficiency and performance, and accordingly, the GDI engines are recently increasingly used.\nVarious methods of mixing air with fuel well and maximally concentrating an air-fuel mixture around the spark plug has been proposed to smoothly operate the engine at a small air/fuel ratio.\nA vortex is generated in the axis perpendicular to the movement direction of the piston or with respect to the movement direction axis of the piston in the internal combustion engine, in which the former one is called ‘tumble’ and the later one is called ‘swirl’.\nSince the mixing ratio and concentration of the air and fuel depends on the flow level of the tumble and the swirl, design should be performed in consideration of the tumble and swirl to improve the operational performance of the GDI engine.\nThe tumble and swirl particularly depends on the shape of the top of the piston, the design of the top of the piston should be improved to improve the operational performance of the GDI engine.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
{"text": "Modern wireless communication networks typically allocate resources to various devices by allocating certain frequencies to the devices at various time slots. However, it is often the case that uplink transmission signals of one or more devices on the network leak into non-allocated frequencies. Such signals can interfere with communications of other devices on the network and with communications of devices on other networks."}
{"text": "1. Field of the Invention\nThe present invention relates to an image processing device and image processing method for computer graphics.\n2. Description of the Related Art\nIn the field of computer graphics, various representations are performed on the screen. One of these representations is the residual-image effect, whereby part of a past image is allowed to remain behind a moving object. Conventionally, residual image processing was performed by software. Residual image processing is performed by gradually altering the transparency of the moving object that is displayed, in accordance with its movement. A number of polygons equal to (number of polygons of the object)xc3x97(number of residual images) are therefore required in order to represent residual images. In the conventional signal processing device, the residual image portion was calculated in the same way as an ordinary object, and rendering had to be carried out, so the efficiency of processing was poor. Also, if a large number of polygons are employed, the amount of processing involved is large, so the processing took a long time.\nFurthermore, while it may be intended to produce the same residual image effect for all polygons, there are also cases in which it is desirable that it should be possible to set a desired residual image effect for each polygon or to effect display with desired polygons emphasised.\nThe present invention was made in order to solve the above problems. An object of the present invention is to provide an image processing device wherein the efficiency of processing can be raised and a speeding-up of processing can be achieved by giving the residual image function to hardware.\nAlso, an object of the present invention is to provide an image processing device whereby a desired residual image effect etc. can be attained for desired polygons.\nAn image processing device according to the present invention comprises a control section that writes generated image data to a first image memory region and a second image memory region, and further comprises a blend circuit that reads image data from said first image memory region and whereby a residual image effect is generated by attenuating the read image data and writing it together with newly generated image data in said second image memory region.\nAn image processing device according to the present invention comprises: a frame buffer that stores image data; a frame buffer interface that receives pixel data from a preceding-stage circuit and that writes the data to said frame buffer and that reads data from said frame buffer; a blend rate register that stores a blend rate when mixing pixel data from said preceding-stage circuit with pixel data of said frame buffer; and a first blend circuit that mixes pixel data from the preceding-stage circuit in accordance with the blend rate of said blend rate register with pixel data from said frame buffer interface and returns the result to the preceding-stage circuit.\nAn image processing device according to the present invention comprises: a frame buffer that stores, in pixel units, image data, blend rate and slope data, which is a parameter for attenuating blend rate; a frame buffer interface that receives, in pixel units, image data from a preceding-stage circuit, blend rate and slope data, and that writes the data to said frame buffer and that reads the data from said frame buffer; a blend circuit that mixes pixel data from a preceding-stage circuit in accordance with the blend rate of said frame buffer with pixel data from said frame buffer interface; a subtraction circuit that subtracts a value corresponding to slope from the blend rate; a comparison circuit that compares the blend rate of said frame buffer with a blend rate from said subtraction circuit; and a selection circuit that selects one or other of the blend rate and slope of said frame buffer and the blend rate and slope from the preceding-stage circuit in accordance with the output of said comparison circuit and outputs the result to the preceding-stage circuit.\nAn image processing device according to the present invention comprises: a frame buffer that stores image data; a blend rate register that stores a blend rate when pixel data from a preceding-stage circuit and pixel data from said frame buffer are mixed; a blend circuit that mixes pixel data from the preceding-stage circuit with pixel data from said frame buffer in accordance with the blend rate of said blend rate register and that returns the result to the preceding-stage circuit; and a selector that receives a flag from the preceding-stage circuit and that selects one of the image data from the preceding-stage circuit or image data output from said blend circuit and that writes the selected image data to said frame buffer.\nAn image processing device according to the present invention comprises: a control section that writes generated image data to a first image memory region and a second image memory region, further comprising: means for writing in the other image memory one of mixed image data obtained by attenuating image data read from one of said first image memory region and second image memory region and mixing it with newly generated image data, and said newly generated image data.\nAn image processing method for applying a residual image effect to an entire screen according to the present invention comprises: setting the density of the residual image, indicating what percentage of the image data of the previous occasion is to be displayed as residual image, and the separation of residual images, which determines whether the image data to be left as residual image is to be constituted by the data of the previous occasion or the data of the occasion previous to the previous occasion; blending current pixel colour and buffer pixel colour in accordance with the set residual image density and writing the result to a frame buffer; determining in respect of all pixels in a region whether processing has been completed or not, and, if processing of all the pixels has not been completed, repeating said writing step; and determining whether processing of all the regions has been completed, and, if processing in respect of all the regions has not been completed, repeating said writing step.\nAn image processing method for applying a residual image effect in polygon units according to the invention comprises: setting the density of a residual image, indicating what percentage of the image data of the previous occasion is to be displayed as residual image and the separation of residual images, that determines whether the image data that is to be left behind as residual image is that of the data of the previous occasion or of the data of the occasion previous to the previous occasion; sorting the polygons contained in the image; blending the current pixel colour and buffer pixel colour in accordance with the set density of the residual image and writing the result to the frame buffer; comparing the density of the image of the current pixel with that when the slope data, which is a parameter indicating the degree of attenuation from the density of the image, is subtracted, and writing the larger image density and slope data in the residual image buffer; determining whether or not processing has been performed in respect of all the pixels in the region and, if processing has not been completed in respect of all pixels, repeating said writing step; and determining whether processing has been completed in respect of all regions, and, if processing has not been completed in respect of all regions, repeating said writing step."}
{"text": "Protein C is a serine protease and naturally occurring anticoagulant that plays a role in the regulation of homeostasis by deactivating Factors V.sub.a and VIII.sub.a in the coagulation cascade. Human protein C is made in vivo primarily in the liver as a single polypeptide of 461 amino acids. This precursor molecule undergoes multiple post-translational modifications including 1) cleavage of a 42 amino acid signal sequence; 2) proteolytic removal from the one chain zymogen of the lysine residue at position 155 and the arginine residue at position 156 to make the 2-chain form of the molecule, (i.e., a light chain of 155 amino acid residues attached through a disulfide bridge to the serine protease-containing heavy chain of 262 amino acid residues); 3) vitamin K-dependent carboxylation of nine glutamic acid residues clustered in the first 42 amino acids of the light chain, resulting in 9 gamma-carboxyglutamic acid residues; and 4) carbohydrate attachment at four sites (one in the light chain and three in the heavy chain). The heavy chain contains the well established serine protease triad of Asp 257, His 211 and Ser 360. Finally, the circulating 2-chain zymogen is activated in vivo by thrombin at a phospholipid surface in the presence of calcium ion. Activation results from removal of a dodecapeptide at the N-terminus of the heavy chain, producing activated protein C (aPC) possessing enzymatic activity.\nIn conjunction with other proteins, protein C functions as perhaps the most important down-regulator of blood coagulation. In other words the protein C enzyme system represents a major physiological mechanism of anticoagulation.\nThe coagulation system is best viewed as a chain reaction involving the sequential activation of zymogens into active serine proteases eventually producing the enzyme, thrombin, which through limited proteolysis converts plasma fibrinogen into the insoluble gel, fibrin. Two key events in the coagulation cascade are the conversion of clotting factor X to Xa by clotting factor IXa and the conversion of prothrombin into thrombin by clotting factor Xa. Both of these reactions occur on cell surfaces, most notably the platelet surface. Both of these reactions require cofactors. The major cofactors, factors V and VIII, in the system circulate as relatively inactive precursors, but when the first few molecules of thrombin are formed, thrombin loops back and activates the cofactors through limited proteolysis. The activated cofactors, Va and VIIIa, accelerate both the conversion of prothrombin into thrombin and also the conversion of factor X to factor Xa by approximately five orders of magnitude. Activated protein C overwhelmingly prefers two plasma protein substrates which it hydrolyzes and irreversibly destroys. These plasma protein substrates are the activated forms of the clotting cofactors, Va and VIIIa. Activated protein C only minimally degrades the inactive precursors, clotting factors V and VIII. Activated protein C in dogs has been shown to sharply increase circulating levels of the major physiological fibrinolytic enzyme, tissue plasminogen activator (tPA). Activated protein C has been shown in vitro to enhance the lysis of fibrin in human whole blood. Therefore, activated protein C represents an important adjunct to in vivo fibrinolysis in man.\nToday, there are few effective treatments available for vascular occlusions, including thrombotic stroke. Treatment with tPA, if administered within three hours from the onset of the stroke, has been recently approved by the FDA. Treatment of strokes with either heparin or oral anticoagulants, although occasionally beneficial, carries a high risk for bleeding into the infarcted brain area.\nThe use of recombinant aPC (r-aPC) in the treatment of thrombotic occlusion or thromboembolism in a baboon model has been presented by Griffin, et al. in U.S. Pat. No. 5,084,274 and European Patent Specification EP 0 318 201 B1. Griffin claimed dose levels in the range of 0.07 mg/kg/hr to 1.1 mg/kg/hr for the treatment of thrombotic occlusion. However, applicants have found that these dose levels are in a range above the toxicological level of r-aPC. For example, pre-clinical toxicology studies in non-human primates indicate the safety of r-aPC for a 96 hour infusion is limited at a top dose of around 0.05 mg/kg/hr. Therefore, the lowest dose level taught by Griffin, et al., i.e. 0.07 mg/kg/hr, is at a level greater than the toxic dose established by applicants for humans. Thus, even the lowest dose level taught by Griffin would carry a high risk for bleeding into the infarcted brain area, thereby aggravating the neurological deficit accompanying the stroke. Accordingly, even in view of the teaching of Griffin, et al., there remains a need to identify an effective therapy of arterial thrombus formation in humans with aPC.\nContrary to the teachings of prior investigators, applicants have discovered that only low dose therapy with r-aPC is useful in the treatment of thrombotic stroke. The administration of aPC is also beneficial in preventing the local extension of the microvascular and macrovascular occluding arterial thrombus, thereby reducing the neurological deficit resulting from the stroke."}
{"text": "Power supplies, such as switching power supplies, are used for providing direct current (DC) supply for various electronic devices, such as base stations in a cellular communications system. Such power supplies have certain output current limits, which must not be exceeded. It is common to use circuit breakers in order to protect power supplies from excessive currents due to overloads. The circuit breakers also protect wiring in high energy circuits and limit outages in case of fault conditions. A fault in one unit thus does not cause breakdown of other units of the same power supply.\nA common type of a circuit breaker is an electromechanical circuit breaker. Such a circuit breaker has a fixed current limit or fixed tripping curve, and exceeding current limit or tripping curve causes the circuit breaker to trip. A circuit breaker usually has a lever at its front panel for resetting the circuit breaker. Electromechanical circuit breakers are produced in large quantities, and they are standard, reliable and low cost components, which have been approved by authorities.\nHowever, there are some problems related to the electromechanical circuit breakers. After tripping they need to be manually reset. Electrical devices may be installed in remote locations, and it may require much effort and time to travel to such locations. Tripping of a circuit breaker may be caused by reasons which do not require repairing the electrical devices. In such situations it is still necessary to travel to the installation location in order to reset circuit breaker. To solve this problem there are arrangements for resetting a circuit breaker with remote control. Such a solution is disclosed in patent application document U.S. Pat. No. 6,522,227 B1.\nAnother problem with electromechanical circuit breakers relates to the fixed tripping current limit or fixed tripping curve. There are various situations where the current limit of a circuit breaker needs to be changed. In order to change the current limit the circuit breaker must be changed to a different one. Changing a circuit breaker requires servicing work at the installation location of the electrical device. It is also possible that a circuit breaker with the required current limit is not available. In such case it is necessary to use a circuit breaker with non-optimal, lower current limit. This may cause unnecessary tripping of the circuit breaker."}
{"text": "1. Field of the Invention\nThe present invention generally relates to laser-induced breakdown spectroscopy systems\n2. Description of Related Art\nLaser-induced breakdown spectroscopy (“LIBS”) is a type of atomic emission spectroscopy which uses a highly energetic laser pulse as the excitation source. The laser is focused to form a plasma, which atomizes and excites samples. In principle, LIBS can analyze any matter regardless of its physical state, be it solid, liquid, or gas. Because all elements emit light of characteristic frequencies when excited to sufficiently high temperatures, LIBS can detect all elements, limited only by the power of the laser beam utilized as well as the sensitivity and wavelength range of the spectrograph and detector.\nIf the constituents of a material to be analyzed are known, LIBS may be used to evaluate the relative abundance of each constituent element, or to monitor the presence of impurities. In practice, detection limits are a function of a) the plasma excitation temperature, b) the light collection window, and c) the line strength of the viewed transition. LIBS makes use of optical emission spectrometry and is to this extent very similar to arc/spark emission spectroscopy.\nLIBS operate by focusing the laser beam onto a small area at the surface of the specimen When the laser beam is discharged it ablates a very small amount of material, in the range of nanograms to picograms, which generates a plasma plume with temperatures in excess of 100,000 K. During data collection, typically after local thermodynamic equilibrium is established, plasma temperatures range from 5,000-20,000 K. At the high temperatures during the early plasma, the ablated material dissociates (breaks down) into excited ionic and atomic species. During this time, the plasma emits a continuum of radiation which does not contain any useful information about the species present, but within a very small timeframe the plasma expands at supersonic velocities and cools. At this point the characteristic atomic emission lines of the elements can be observed."}
{"text": "1. Field of the Invention\nThis invention relates to mechanized equipment for assisting in making pleated window draperies by mass production, and particularly to a machine for marking the material required by the pleats and for marking the material required for the spaces between the pleats along the top edge of a drapery panel and for simultaneously marking the hem of the pleated panel to assist in folding the drape.\n2. Description of the Prior Art\nThe most common manner of making draperies employs the hand-measuring of the various panels of drapery material. In recent years, precision machines have been devised to speed up the various operation, both to save labor and reduce the need for skilled labor.\nThe Caracoes U.S. Pat. No. 3,145,480 discloses a machine for marking the material required to form pleats along the top edge of a drapery panel. Caracoes accomplishes this by a series of pointers mounted on a single lazy tong apparatus. This machine marks (by pointers) the space required for both the pleat and the space between the pleat.\nThe Tuskos U.S. Pat. No. 3,696,515 shows a drapery pleats marking machine with a workbench that supports two vertically-disposed, lazy tong mechanisms. A marking pen is assembled at each intermediate pivot point of both lazy tong mechanisms. The adjustment of the first lazy tong mechanism changes the spacing P allotted for making material into the pleats, while the adjustment of the second lazy tong mechanism changes the spacing S. for the material required for forming the space between the pleats."}
{"text": "1. Field of the Invention\nThis invention relates to self-regulating pressure systems and methods for vehicle tires, and more particularly to self-regulating systems and methods in which the tire is automatically replenished from a high pressure reservoir in response to a sensed low pressure condition in the tire.\n2. Description of the Related Art\nLow tire pressure is an important cause of excessive fuel consumption, tire wear and impaired steerability. A normal tire will typically leak on the order of 25 percent of its pressure per year due to its inherent permeability. It is thus good practice to check tire pressure on a regular basis. However, even checking tire pressure every few weeks may not prevent these adverse affects when a slow leak is present, and the leak may go undetected unless a careful record is maintained of how frequently the pressure in each tire has to be replenished. A fast leak or flat condition can rapidly cause damage to the tire and even render it unusable in a short period of time, but this condition may go unnoticed by an inexperienced driver until it is too late.\nIt is thus highly desirable to have some mechanism that automatically replenishes the tire pressure when it is too low, and that warns the driver of the low pressure condition. One such system is disclosed in U.S. Pat. No. 4,067,376 to Barabino. This patent incorporates a high pressure reservoir into the vehicle wheel, and uses a valve that automatically opens a passageway between the high pressure reservoir and the tire in response to the tire pressure falling below a selected threshold level. Operation of the valve also causes a sonic or ultrasonic signal to be generated which is detected by a sensor mounted in close proximity to the wheel, typically within the wheel well. When it senses a valve operation, the sensor initiates a visual and audio warning for the driver. To differentiate between the different tires, the valve on each tire may be tuned to generate a signal with a unique frequency, or the sensors for the different tires may be tuned to detect different components of a uniform valve signal.\nAnother system is disclosed in U.S. Pat. No. 4,742,857 to Gandhi. In this patent the tire valve includes a magnet that moves in response to a falling tire pressure. Movement of the magnet to a low pressure position is sensed by a solenoid mounted on the shock absorber adjacent the tire. The magnetic field from the moveable magnet intercepts the solenoid at each tire revolution, generating a voltage in the solenoid that varies with the magnet's position. When the magnet's sensed position indicates a low pressure condition, a warning is provided to the vehicle operator. If the operator wishes to replenish the tire pressure, he or she actuates a switch that supplies current to the solenoid for the appropriate tire. The solenoid produces a magnetic field that moves its associated magnet to a position at which a valve from a high pressure reservoir within the wheel is opened so that air can flow into the tire.\nWhile the above systems can effectively replenish the air pressure within a tire, they are incapable of distinguishing between normal long term leakage from a tire, a slow leak condition that requires repair, and a fast leak or flat tire condition. The vehicle operator is thus uninformed as to whether a tire repair or replacement is called for. The operator is also not supplied with any information as to the condition of the high pressure reservoir, and whether it requires replenishment. Furthermore, the prior use of separate transmitters within each wheel and associated receivers on the vehicle requires a power source within each wheel and complex electronics, and is unduly expensive.\nAn after market sensor for low tire pressure has been marketed by Epic Technologies Inc. The sensor monitors tire pressure and causes a radio frequency (RF) signal to be transmitted in case the pressure falls below a predetermined set point; transmission of the RF signal actuates a warning for the vehicle driver. However, the system does not have any mechanism to automatically replenish the low tire pressure."}
{"text": "Modern personal computers (PC) are increasingly used in companies, offices and by private users. Their failure or malfunction may cause severe and costly damage or expensive loss of data and information. Unfortunately, the quality of utility power does not guaranty reliable operation of personal computers. Outages, abnormal voltage and transients disturb or halt the operation, erase precious data and cause deterioration of components and hardware failures."}
{"text": "A remote operator station (ROS) for a machine allows an operator to control the machine from a remote location. The machine may be employed on an actual worksite while the ROS may be positioned indoors. U.S. Pat. No. 6,633,800 discloses an apparatus and method for a control unit which allows for autonomous, manual and tele-operation of mining vehicles. The control unit has a robust system design to withstand the harsh environment of underground mines The control unit allows a tele-operator, in a remote tele-operator station, to use image and operational data, joysticks and foot pedals to remotely control the mining vehicle. In another aspect, the control unit provides safety features such as supervising its operation for operational errors and providing status, warning and error information to the tele-operator station."}
{"text": "DE 10 2007 037 204 A1 discloses a process for manufacturing an optical lens element, in particular for illumination purposes, particularly for manufacturing a headlight lens for a vehicle headlight, in particular for a motor vehicle headlight, wherein a blank is molded from transparent, particularly thermoplastic, in particular essentially liquid plastic material, in particular by means of an injection pressing process in an injection pressing mold, and wherein the blank is subsequently pressed, press-molded by means of a final shape/contour mold to form the lens element.\nDE 41 28 915 A1 discloses an apparatus for casting a plastic lens, wherein a first casting mold having a concave optical surface, and a second casting mold having a convex optical surface are arranged such that the optical surfaces oppose each other, and wherein an adhesive tape is placed over the first as well as also over the second casting mold in order to configure a cavity between the casting molds. Plastic material is injected into this cavity and is finally hardened.\nUS 2006/0284324 A1 discloses a method for producing an optical structural component wherein the structural component is pressed after thermal softening.\nU.S. Pat. No. 3,211,811 discloses an apparatus for producing a plastic-material lens, wherein a flexible molding element is arranged in a mold cavity, which element divides the mold cavity into two partial mold cavities, wherein the flexible mold element is configured such that it transmits a hydraulic force between the partial mold cavities.\nDE 699 23 847 T2 discloses a process for manufacturing an optical mold body from thermoplastic resin by injection-press-molding, wherein the volume of a mold cavity expands more strongly than the volume of the optical mold body in question, and molten thermoplastic resin is injected into the mold cavity through an injection molding passage.\nAccording to DE 699 23 847 T2, U.S. Pat. No. 4,540,534, EP 0 640 460 and JP 9-057794 disclose a method for manufacturing an optical mold body from thermoplastic resin by injection-press-molding.\nDE 102 20 671 A1 discloses a plastic lens consisting of a collective lens of plastic material having a high Abbé number, and a dispersion lens connected therewith integrally and in positive fit, of plastic material having, when compared with the plastic material of the collective lens, a lower Abbé number, wherein the thermal expansion coefficients of the respective plastic materials are essentially the same.\nHeadlight lenses are known e.g. from WO 02/31543 A1, U.S. Pat. No. 6,992,804 B2, WO 03/074251 A1 and DE 100 52 653 A1. Further vehicle headlights are known e.g. from DE 100 33 766 A1, EP 0 272 646 A1, DE 101 18 687 A1 and DE 198 29 586 A1.\nIt is an object of the invention to reduce the costs for manufacturing of headlight lenses for vehicle headlights. It is, for example, an object of the invention to specify an improved vehicle headlight including a plastic lens and/or a motor vehicle with a correspondingly improved vehicle headlight."}
{"text": "Countless products are packaged and shipped to end-users in this country and around the world daily. Many products are placed in crates or boxes filled with packing material to minimize or prevent damage during shipping. Some products are stored and packaged in cartons, which may be sealed to prevent the enclosed items from exposure to ambient conditions. Other products are less sensitive and don't require a hermetic seal. Examples of packaged products may range from edible substances to industrial consumables, like welding electrodes. Cartons are commonly constructed from renewable or recycled material like, for example, corrugated cardboard, which is easy and relatively inexpensive to produce.\nIn many cases, similarly sized cartons are used to package different quantities of products. Some cartons may therefore contain a particular quantity of material leaving free space within the carton for the products to jostle about during shipment resulting in damage. This also affords the opportunity for stacked cartons to compress or deform. As the contents may not be uniformly distributed, pressure frequently collapses one end of the adjacent carton creating an uneven stacking surface. Cartons higher in the stack are therefore prone to falling and hence damage. Falling cartons also lead to an unsafe working environment. Furthermore, the tilted packaged contents look disorderly and unsightly.\nIt would be useful to provide a carton that automatically collapses under pressure in a uniform manner. Such a container would limit spacing within the container and maintain a level relationship between upper and lower surfaces of adjacent cartons."}
{"text": "1. Field of the Invention\nImplementations described herein relate generally to network security and, more particularly, to security perimeters that may be used to protect a customer network from network attacks.\n2. Description of Related Art\nNetwork attacks are a real, growing threat to businesses worldwide. These attacks can quickly incapacitate a targeted business, costing victims thousands, if not millions, of dollars in lost revenue and productivity.\nOne type of network attack, called a denial of service attack, can paralyze Internet systems by overwhelming servers, network links, and network devices (e.g., routers, firewalls, etc.) with bogus traffic. Easily launched against limited defenses, network attacks not only target individual web sites or other servers at the edge of the network, but they also can incapacitate the network itself.\nThe growing dependence on the Internet makes the impact of successful network attacks increasingly painful (financially and otherwise) for service providers, enterprises, and government agencies. Newer, more powerful tools used by network attackers promise to unleash even more destructive attacks in the months and years to come.\nBecause denial of service attacks are among the most difficult to defend against, responding to them appropriately and effectively poses a tremendous challenge for all Internet-dependent organizations. Network devices and current perimeter security technologies, such as firewalls and intrusion detection devices (IDDs), although important components of an overall security strategy, do not by themselves provide comprehensive network attack protection.\nCurrent techniques fall short in terms of prevention and/or mitigation of network attacks. Some of the more popular network attack responses, such as blackholing and router filtering, are not optimized to deal with increasingly sophisticated attacks. IDDs offer some excellent attack detection capabilities, but cannot mitigate the impact of the attacks. Firewalls offer a rudimentary level of protection but, like blackholing and router filtering, they were not designed to protect against the types of advanced attacks that are so common today. Still other strategies, such as overprovisioning, do not provide adequate protection against ever larger attacks, and can be costly as a network attack prevention strategy.\nNetwork attacks also make it difficult to keep up with the constant system patching and upgrading required for a successful security system. Many times, implementing significant patches or upgrades is a disruptive event within a customer's network. While overall service disruption can be avoided in a multiple server environment, it still takes time to determine the impact of the patch on the customer's applications and then perform the upgrades."}
{"text": "The use of liposomes for drug delivery has been proposed for a variety of drugs, particularly those which are administered parenterally. Liposomes have the potential for providing controlled \"depot\" release of the administered drug over an extended time period, and of reducing side effects of the drug, by limiting the concentration of free drug in the bloodstream. Liposomes can also alter the tissue distribution and uptake of drugs, in a therapeutically favorable way, and can increase the convenience of therapy, by allowing less frequent drug administration. Liposome drug delivery systems are reviewed in Poznansky.\nGenerally, the optimal liposome size for use in parenteral administration is between about 0.1 and 0.3, and up to 0.4, microns. Liposomes in this size range can be sterilized by passage through conventional depth filters having particle size discrimination of about 0.2 microns. This size range of liposomes also favors biodistribution in certain target organs, such as liver, spleen, and bone marrow (Gabizon), and gives more uniform and predictable drug-release rates and stability in the bloodstream. Liposomes whose sizes are less than about 0.4 microns also show less tendency to aggregate on storage, and are thus generally safer and less toxic in parenteral use than larger-size liposomes.\nA variety of techniques have been proposed for preparing liposomes, including drug-containing liposomes (Szoka, 1980). Typically, these methods yield liposomes which are heterodisperse, and predominantly greater than about 1 micron in size. These initial heterodisperse suspensions can be reduced in size and size distribution by a number of known methods. One size-processing method which is suitable for large-scale production is homogenization. Here the initial heterodisperse liposome preparation is pumped under high pressure through a small orifice or reaction chamber. The suspension is usually cycled through the reaction chamber until a desired average size of liposome particles is achieved. A limitation of this method is that the liposome size distribution is typically quite broad and variable, depending on a number of homogenization cycles, pressures, and internal temperature. Also, the processed fluid has the potential to pick up metal and oil contaminants from the homogenizer pump, and may be further contaminated by residual chemical agents used to sterilize the pump seals.\nSonication, or ultrasonic irradiation, is another method that is used for reducing liposome sizes. This technique is useful especially for preparing small unilameller vesicles (SUVs), in the 0.025-0.08 micron size range. However, a narrow size distribution of liposomes can only be achieved at liposome sizes of about 0.05 microns, i.e., after the liposomes have been reduced to their smallest sizes. The very small liposomes have limited drug carrying or loading capacity and less favorable biodistribution properties than those in the 0.1-0.4 micron size range, as noted below. The processing capacity of this method is also quite limited, since long-term sonication of relatively small volumes is required. Also, heat build-up during sonication can lead to peroxidative damage to lipids, and sonic probes shed titanium particles which are potentially quite toxic in vivo.\nA third general size-processing method known in the prior art is based on liposome extrusion through uniform pore-size polycarbonate membranes (Szoka, 1978). This procedure has advantages over the above homogenization and sonication methods in that a variety of membrane pore sizes is available for producing liposomes in different selected size ranges, and in addition, the size distribution of the liposomes can be made quite narrow, particularly by cycling the material through the selected-size filter several times. Nonetheless, the membrane extrusion method has several drawbacks in large-scale processing. For one, the pores in the membrane tend to clog, particularly when processing concentrated suspensions and/or when the liposome sizes are substantially greater than the membrane pore sizes. The clogged membranes cannot be cleared, because the filter housing configuration does not allow back flushing. Replacing the filter is likely to compromise the sterility of the extrusion operation and the membranes cannot be steam-sterilized in place, with a high degree of confidence, due to their inherent fragility. Secondly, the total pore surface area is only about 20% of the total surface area of the membrane, restricting the surface area available for extrusion, and thus limiting total throughput.\nCo-owned U.S. Pat. No. 4,737,323 to \"Liposome Extrusion Method\" describes a method for sizing liposomes by extrusion through an asymmetric ceramic filter. According to one aspect of the invention, a single pass through the filter converts a suspension of heterogeneous size liposomes with a substantial portion having sizes greater than 1 micron to a relatively homogeneous population having sizes less than about 0.4 micron. One or more additional passes through the filter, especially in an inside-to-outside direction, produce a convergence to a final reduced average size and polydispersity. According to another aspect of the invention, it was found that the average size of the extended liposomes was substantially smaller than the rated pore size of the ceramic filter.\nCeramic filter extrusion is advantageous for high-volume processing, since the filter presents a large surface area, and the pore surface area constitutes a major portion of total area. Further the filter is designed for operation at relatively high pressure. In addition, the filter can be backflushed to prevent clogging, and easily sterilized. The method thus overcomes the major limitations of the polycarbonate membrane extrusion method, while producing comparable liposome sizing."}
{"text": "The advantage of the approach taken by the present invention is to produce single immunogens containing protective antigens from a range of pathogens. Such chimeras greatly simplify the development of combination vaccines, in particular, with the view ultimately to produce single dose multivalent vaccines. Multivalent vaccines are currently made by separately producing pathogens and/or their pertinent antigens and combining them in various formulations. This is a labour intensive, costly and complex manufacturing procedure. In contrast, the availability of a single immunogen capable of protecting against a range of diseases would solve many of the problems of multivalent vaccine production. Several chimeric immunogens of the type provided herein may be combined to decrease the number of individual antigens required in a multivalent vaccine.\nHuman Parainfluenza virus types 1,2,3 and Respiratory syncytial virus types A and B are the major viral pathogens responsible for causing severe respiratory tract infections in infants and young children. It is estimated that, in the United States alone, approximately 1.6 million infants under one year of age will have a clinically significant RSV infection each year and an additional 1.4 million infants will be infected with PIV-3. Approximately 4000 infants less than one year of age in the United States die each year from complications arising from severe respiratory tract disease caused by infection with RSV and PIV-3. The WHO and NIALD vaccine advisory committees ranked RSV number two behind HIV for vaccine development while the preparation of an efficacious PIV-3 vaccine is ranked in the top ten vaccines considered a priority for vaccine development.\nSafe and effective vaccines for protecting infants against these viral infections are not available and are urgently required. Clinical trials have shown that formaldehyde-inactivated and live-attenuated viral vaccines failed to adequately protect vaccines against these infections. In fact, infants who received the formalin-inactivated RSV vaccine developed more serious lower respiratory tract disease during subsequent natural RSV infection than did the control group. [Am. J. Epidemiology 89, 1969, p.405-421; J. Inf. Dis. 145, 1982, p.311-319]. Furthermore, RSV glycoproteins purified by immunoaffinity chromatography using elution at acid pH induced immunopotentiation in cotton rats. [Vaccine, 10(7), 1992, p.475-484]. The development of efficacious PIV-3 and RSV vaccines which do not cause exacerbated pulmonary disease in vaccines following injection with wild-type virus would have significant therapeutic implications. It is anticipated that the development of a single recombinant immunogen capable of simultaneously protecting infants against diseases caused by infection with both Parainfluenza and Respiratory syncytial viruses could significantly reduce the morbidity and mortality caused by these viral infections.\nIt has been reported that a protective response against PIV-3 and RSV is contingent on the induction of neutralizing antibodies against the major viral surface glycoproteins. For PIV, these protective immunogens are the HN protein which has a molecular weight of 72 kDa and possesses both hemagglutination and neuraminidase activities and the fusion (F) protein, which has a molecular weight of 65 kDa and which is responsible for both fusion of the virus to the host cell membrane and cell-to-cell spread of the virus. For RSV, the two major immunogenic proteins are the 80 to 90 kDa G glycoprotein and the 70 kDa fusion (F) protein. The G and F proteins are thought to be functionally analogous to the PIV HN and F proteins, respectively. The PIV and RSV F glycoproteins are synthesized as inactive precursors (FO) which are proteolytically cleaved into N-terminal F2 and C-terminal F1 fragments which remain linked by disulphide bonds.\nRecombinant surface glycoproteins from PIV and RSV have been individually expressed in insect cells using the baculovirus system [Ray et al., (1989), Virus Research, 12: 169-180; Coelingh et al., (1987), Virology, 160: 465-472; Wathen et al., (1989), J. of Inf. Dis. 159: 253-263] as well as in mammalian cells infected with recombinant poxviruses [Spriggs, et al., (1987), J. Virol. 61: 3416-3423; Stott et al., (1987), J. Virol. 61: 3855-3861]. Recombinant antigens produced in these systems were found to protect immunized cotton rats against live virus challenge. More recently, hybrid RSV F-G [Wathan et al., (1989), J. Gen Virol. 70: 2625-2635; Wathen, published International Patent application WO 89/05823] and PIV-3 F-HN [Wathen, published International Patent Application WO 89/10405], recombinant antigens have been engineered and produced in mammalian and insect cells. The RSV F-G hybrid antigen was shown to be protective in cotton rats [Wathan et al., (1989), J. Gen. Virol. 70: 2637-2644] although it elicited a poor anti-G antibody response [Connors et al., (1992), Vaccine 10: 475-484]. The protective ability of the PIV-3 F-HN protein was not reported in the published patent application. These antigens were engineered with the aim to protect against only the homologous virus, that is either RSV or PIV-3. However, it would be advantageous and economical to engineer and produce a single recombinant immunogen containing at least one protective antigen from each virus in order simultaneously to protect infants and young children against both PIV and RSV infections. The chimeric proteins provided herein for such purpose also may be administered to pregnant women or women of child bearing age to stimulate maternal antibodies to both PIV and RSV. In addition, the vaccine also may be administered to other susceptible individuals, such as the elderly."}
{"text": "The present invention relates to a scanning tunnel microscope and more particularly to a scanning tunnel microscope which is designed in such a way that a position of a desired portion on a specimen to be observed with the scanning tunnel microscope is previously positioned by utilizing a scanning electron microscope so that an observation can be then made with respect to the positioned portion.\nWith respect to a prior art apparatus which is constructed by combining a scanning tunnel microscope (hereinafter, referred to as simply \"an STM\", when applicable) with a scanning electron microscope (hereinafter, referred to as simply \"a SEM\", when applicable), the description is shown in an article \"Review of Scientific Instruments 60\", April, 1989, pp. 789 to 791. In the apparatus described in this article, a specimen is mounted on a specimen holder in such a way that its surface makes an angle of 45.degree. with an electron beam of the SEM. The specimen is then moved towards a probe of the STM by a coarse movement mechanism referred to as an inchworm utilizing a piezoactuator, which serves to linearly move a specimen in a direction perpendicular to the electron beam, so as to approach to a position at a predetermined distance from the probe which allows the tunnel effect to be generated. The probe of the STM is attached to a fine movement mechanism referred to as a piezotripod, so that it is moved along the surface of the specimen, thus performing a scanning operation. The above-mentioned elements constituting the STM are fitted to a stage of the SEM within a specimen chamber of the SEM.\nIn the prior art apparatus described above, there is provided therein only the coarse movement mechanism which serves lineally to make the specimen approach to the probe to be located in an area allowing generation of the tunnel effect, and also a SEM stage for searching for the probe by utilizing the SEM is not clear in construction.\nJP-A-63-116349 (the term JP-A used herein means that the patent application was laid open to public inspection but has not been examined) and JP-A-63-298951 disclose that a position on a specimen to be observed with a STM is previously positioned by utilizing a SEM and that a specimen is adapted to be moved in a two dimensional manner along its surface with its surface being inclined at 45.degree. with an electron beam.\nMoreover, JP-A-64-79603 discloses a microscope which is designed in such a way that a specimen is observed at a low magnification with a SEM and the same specimen is observed at a high magnification with a STM.\nThe prior art described above discloses that a position to be observed with an STM with a higher magnification is determined using a SEM, and then an observation is made with the STM. However, when the apparatus is actually operated, even if the position of the specimen is set in such a way that an observation position determined by the SEM goes to a predetermined position within a visual field of an electron microscope image by the SEM, it is not necessarily easy to make the probe of the STM scan accurately over the observation position determined by the SEM since the position of the probe of the STM is not precisely defined with respect to an image by the SEM.\nIn other words, only when the relationship among the specimen, the probe and the electron beam is accurately determined, the position positioned by the SEM can be observed accurately at a high magnification with the STM."}
{"text": "In order to stabilize a downhole emplacement such as a bored or drilled hole forming a well, it typically requires multiple materials that are dry blended in a pneumatic bulk plant and then injected down hole in a cementing process. These dry materials can include cement, pozzolanic fly ash, silica, chemical additives that alter the cure time (accelerators/retarders), fluid loss materials, viscosifiers, defoamers, friction reducers, extenders, and other chemical additives that alter a product's performance downhole. These materials are blended in the pneumatic bulk plant and require multiple transfers to obtain a product that is sufficiently homogenous. Once the blending process is complete, the performance characteristics of the product are set since the blended mix cannot be re-mixed once it is removed from the bulk plant location.\nThe blended dry product is then transported to the site of the oil/gas well using bulk storage units (e.g. bulk trailers, or transports). The product is then introduced to a pumping and mixing truck, such as by pneumatic conveying. In the mixing truck, water is added to the dry product in specific measured quantities to achieve a desired density product slurry that can be measured using a densitometer. The mixing of water with the dry product is done “on the fly”—that is, continually mixed and pumped: pumping operations are not stopped during mixing operations in which the desired product slurry density may be achieved automatically through an electronic control system. The mixing truck may include controls such as water and cement throttling valves that are adjusted dynamically via the electronic control system to ensure that the measured density of the fluid is in agreement with a desired density that is typically determined by the user.\nWhen the slurry mix has achieved the desired density, it can be pumped downhole within a well form casing that is pre-positioned in the well. A drilling or work over-rig may be used to pump the slurry mix in which an annular space between the wellbore and casing are filled incrementally as the slurry mix is pumped."}
{"text": "The present invention relates to a steam iron having a low temperature pressing surface. The common steam iron of today's market utilizes an aluminum soleplate which provides the pressing surface and also includes a boiler, generally of the flash boiler type. It is necessary to heat this aluminum mass to a sufficiently high temperature, generally over 300.degree. F., when the iron is used as a steam iron in order to assure that the water that enters the flash boiler reaches a sufficiently high temperature to be turned into steam and maintain this state so that it does not leave the iron in water droplets and thus water stain the clothes. The pressing surface which is also a part of this aluminum mass is thus maintained at a relatively high temperature which would be damaging to the low temperature fabrics of today's market. Therefore, in irons of this type steam ironing of low temperature fabrics is not possible and thus the iron must be used as a dry iron when low temperature fabrics are ironed.\nPrior art has also attempted to solve this problem by utilizing a plastic add-on shoe disposed beneath the plate of the iron. Such shoes, however, have not found general acceptance. Furthermore, the irons upon which they are used are not designed with this intent in mind and thus are not calibrated with respect thereto thus making the temperature settings for use with the shoe purely guesswork."}
{"text": "1. Field of the Invention\nThe invention is directed to insect repellent formulations which contain polar solvents, have low viscosity, and are deliverable by piezoelectric device.\n2. Description of the Related Art\nSprayable insecticidal compositions are known and generally comprise an insecticide and a solvent, in addition to other components. As the common insecticides are non-polar, the most common solvents used for this purpose are themselves non-polar, e.g. hydrocarbon solvents, such as hexane and petroleum solvent blends available under the trade names Shellsol and Isopar, for example.\nWhile hydrocarbon solvents are effective to dissolve most common non-polar insecticides, they tend to have low flash points and to otherwise pose environmental and health risks. Therefore, it is desired to develop sprayable insecticidal compositions which avoid or minimize the use of these solvents. An insecticide formulation deliverable by piezoelectric delivery system which minimizes or avoids the use of non-polar solvents would be an advancement in the art.\nLow-flash point solvents typically have high volatility. High volatility results in unintended evaporation of the insecticidal formulationxe2x80x94so-called xe2x80x9clatent losses.xe2x80x9d The problem of latent losses is particularly acute when the formulations are intended for delivery by a piezoelectric device over an extended period of time. For this and other reasons, there is a need in the art for insecticidal formulations using low-volatility solvents which are deliverable by piezoelectric device.\nU.S. Pat. No. Re. 31,927 discloses electrostatically sprayable insecticidal formulations. High-boiling, high flash point hydrocarbon solvents, such as Isopar L are used to dissolve these insecticides, but lesser amounts of a polar component, such as n-butanol or cyclohexanone are added to raise the viscosity and reduce the resistivity of the final formulations. The viscosity of these solutions is believed to be too high for effective delivery using a piezoelectric delivery system.\nU.S. Pat. No. 5,935,554 is directed to an aerosol space spray which, in some embodiments, is an insecticidal composition. The composition includes an insecticide, a solvent and a propellant, and in some cases a co-solvent. Where an insecticide is used, the solvent is Shellsol T, with ethanol as a co-solvent.\nEuropean Patent No. 897,755 discloses aerosol and piezoelectric-deliverable insecticide compositions. Mentioned as solvents for the piezoelectric systems are aliphatic hydrocarbons, alicyclic hydrocarbons, aromatic hydrocarbons, halogenated hydrocarbons, alcohols, esters, ethers, and ketones, with a particular preference given to aliphatic saturated hydrocarbons having 5 to 18 carbon atoms. EP 897,755 does not identify polar solvents that would dissolve insecticides, nor does it identify the properties that render a polar solvent suitable for use with a piezoelectric delivery system.\nTypical prior art liquid atomizers utilize a large amount of propellants in the formulation. Conventionally, to obtain the desired characteristics of an aerosol for dispersion into the air, the formulations contain minimal active ingredients, a solvent for the active ingredient and a suitable propellant. The active ingredient usually constitutes less than 1% w/w. The solvent is present in the order of 10-20% w/w and the propellant constitutes 80-90% w/w. In most cases of such aerosol formulations, the solvent is a hydrocarbon solvent and the propellant a fluorocarbon or hydrocarbon. Alternatively, the propellant is partially substituted with water, wherein the weight percentage being water is in the range of 30-40% w/w. Therefore it is evident that to deliver one gram of insecticide conventionally requires the release to the atmosphere of between 80 and 250 g of volatile organic compounds (VOC\"\"s). Volatile organic compound (VOC) is the general name given for compounds with an appreciable vapor pressure, for example, fluorocarbons, hydrocarbons, e.g. butane, etc., commonly used as propellants.\nIt is an object of the invention to provide an insecticidal composition which can be delivered with a piezoelectric delivery system. In particular, it is desired to provide an insecticidal composition with a sufficiently low viscosity that it can be delivered with a piezoelectric delivery system in the form of a fine mist or puff to provide continuous insect repellency throughout a room without settling too rapidly on surrounding surfaces.\nAnother object of the invention is to provide a low viscosity insecticidal composition that is not harmful or hazardous and which minimizes or eliminates the use of hydrocarbon solvents. In particular, it is desired to provide an insecticidal composition that exhibits overall low toxicity to humans and which will not be harmful to the respiratory system. It is also an object of the invention to provide an insecticidal composition that is non-flammable, and exhibits sufficiently low vapor pressure that latent losses are minimized when the formulation is used with a piezoelectric delivery device.\nThese and other objects of the invention are achieved by providing an insecticidal composition comprising an insecticide and at least one polar solvent having a vapor pressure below about 5 mm Hg at 25xc2x0 C., where the insecticidal composition has a viscosity lower than 2.6 centipoise at 25xc2x0 C. In preferred embodiments, the polar solvent comprises dipropylene glycol dimethyl ether.\nPreferably, more than 75% of the solvent is a polar solvent or combination of polar solvents and co-solvents. In particularly preferred embodiments, the composition is substantially free of hydrocarbon solvents.\nIn preferred embodiments, the insecticidal composition according to the invention comprises an insecticide in an amount between about 3 percent by weight and about 20 percent by weight of the composition, while the solvent or solvents comprise between about 80 and 97 percent by weight of the composition.\nThe preferred formulations according to the invention are not emulsions and do not require propellants for delivery.\nAs understood by those of skill in the art, a piezoelectric liquid atomizer comprises an orifice plate which is vibrated by an actuator made of a piezoelectric material, such as lead zirconate titanate (PZT) or lead metaniobate (PN). Upon application of an alternating electrical field across the piezoelectric actuator it expands and contracts and causes the orifice plate to vibrate. The vibrating orifice plate contacts a liquid and as a result liquid droplets are propelled ultrasonically through the orifices in the plate and ejected into the atmosphere.\nThe diameter of the orifices in the orifice plate of a piezoelectric liquid delivery system determines the size of the droplets produced. The diameter generally ranges between about 1 and about 25 microns, preferably from about 4 to about 10 microns. The size of the orifices is tailored such that energy consumption by the device is minimal while the size of the droplets is small enough to permit the droplets to remain suspended in the air and evaporate without collecting on the surrounding surfaces.\nDroplets having an average size greater than about 15 to 20 microns and larger have a tendency to settle out of the air, onto the delivery device and other surfaces. Thus, it is advantageous to provide a solution that will produce droplets having a size less than 20 microns, preferably less than 15 microns. An average droplet size of about 10 microns has been observed to stay suspended in air for a suitable length of time to effectively kill insects in a room.\nThe low viscosity insecticidal compositions according to the invention are particularly suitable for piezoelectric delivery systems powered by low voltage batteries, such as conventional xe2x80x9cAxe2x80x9d, xe2x80x9cAAxe2x80x9d, xe2x80x9cAAAxe2x80x9d, xe2x80x9cCxe2x80x9d, xe2x80x9cDxe2x80x9d cell batteries and the like. In a particularly preferred embodiment, the insecticidal composition of the present invention is adapted for use with a piezoelectric delivery system powered by a single 1.5 Volt xe2x80x9cAAxe2x80x9d size battery.\nPiezoelectric systems that may be used to deliver the insecticidal compositions of the present invention are disclosed in co-pending application Ser. Nos. 09/449,601, 09/498,859, 09/518,882, and 09/519,603, all of which are incorporated by reference herein in their entirety. However, the compositions described and claimed herein are not limited to those deliverable by the systems described in the above-referenced patent applications.\nBy way of description, without limiting the invention, the insecticidal composition may be delivered to a room via a piezoelectric delivery system at a rate between about 20 and about 40 mg/hour, with the insecticidally active component making up between about 3 and 20 percent of that amount.\nIt is generally preferred to obtain a delivery rate of about 1 to about 2 mg per hour of active insecticide so that a 10 ml bottle operated 8 hours a day will have a useful life between about 30 and about 70 days.\nThe use of a piezoelectric delivery device capable of prolonged delivery poses constraints on the solvent systems that can be used in insecticidal formulations. These constraints include: (1) the ability to solubilize insecticide, (2) low viscosity, (3) low vapor pressure, (4) high flash point, and (5) low toxicity (generally evidenced by high solubility in water).\nxe2x80x9cSolvent systemxe2x80x9d as used herein means the combination of all the solvents used in the insecticidal composition. xe2x80x9cMain solventxe2x80x9d as used herein means a solvent comprising more than 50 percent by weight of the final formulation. xe2x80x9cCo-solventxe2x80x9d means any solvent used in addition to the main solvent.\nGenerally, the active must be soluble in the solvent system of the formulation in an amount between about 2 percent by weight to about 10 percent by weight to be deliverable by a piezoelectric device.\nIt has been found that for the effective delivery of insecticide by a piezoelectric delivery system, low viscosity is critical. It is believed that a viscosity lower than 10 centipoise (cp) is necessary to obtain a suitable puff or mist using the preferred delivery systems according to the invention. The insecticidal compositions according to the invention preferably exhibit a viscosity of lower than about 2.6 cp, more preferably, lower than about 2.2 cp.\nAs the insecticides are generally high viscosity substances it is necessary to utilize a solvent system having a viscosity at least less than about 2.6 cp, and preferably less than 2.2 cp or lower, in the insecticidal formulations according to the invention.\nIn preparing the compositions of the present invention, the inventors have sought to lessen or eliminate the use of solvents which are harmful to inhale, such as those designated R20, R21 or R22; solvents irritating to the eyes or skin such as those designated R36 and R38; flammable solvents, such as those designated R10; and solvents harmful to the respiratory system when ingested, vomited and subsequently aspirated into the lungs, such as those designated R65. In general, it is desired that solvents having designation R20, R21 or R22 not exceed 25% by weight of the final formula; solvents having designation R36, R38 not exceed 20% of the final formula; and solvents having designation R65 not exceed 10% of the final formula. The R10 designation refers to low-flashpoint materials. It is desired that the formulations according to the invention have less than 20% by weight of solvents designated R10.\nAs used herein, a xe2x80x9clow toxicityxe2x80x9d composition is one containing less than 25 weight percent of R20, R21, or R22 solvents, less than 20 weight percent of R36 or R38 solvents, or less than 10 weight percent R65 solvent, and which has less than 25 weight percent hydrocarbon solvents overall. Limiting the use of solvents having the above-described hazard designations avoids conspicuous label designations, which is important in in formulating a commercially acceptable insecticidal product. Generally, polar solvents having high solubility in water are low-toxicity solvents.\nThe most preferred compositions according to the invention are substantially free of hydrocarbon solvents. xe2x80x9cSubstantially freexe2x80x9d in this context allows for the presence of minor amounts of hydrocarbon solvents such as will not change the fundamental characteristics of the composition.\nHigh solubility in water is a good indicator of low overall toxicity. Therefore, glycol ethers are good candidates to use as solvents or co-solvents in connection with the invention. For example, propylene glycol ethyl ether, propylene glycol n-propyl ether, ethylene glycol ethyl ether, and ethylene glycol propyl ether all have viscosities below 2.6 centipoise at 25xc2x0 C., good solubility in water, and vapor pressure below 5 mm Hg at 20xc2x0 C. Propylene glycol methyl ether has good water solubility and low viscosity; although its high vapor pressure (8.1 mm Hg at 20xc2x0 C.) precludes its use as a main solvent, it can be used as a co-solvent. Ethylene glycol butyl ether, dipropylene glycol methyl ether, dipropylene glycol n-propyl ether, dipropylene glycol t-butyl ether, and tripropylene glycol methyl ether can also be used, although the viscosities of these glycol ether solvents are all greater than 2.6 cp, which may limit their use somewhat.\nDiethylene glycol methyl ether, diethylene glycol ethyl ether, diethylene glycol propyl ether, and diethylene glycol butyl ether are all water soluble, have low vapor pressures and high flash points. However, their viscosities at 25xc2x0 C. range between 3.9 and 4.7 cp, which may limit their use in formulations according to the invention.\nIt is desirable to keep the flash point of the final formula higher than about 55xc2x0 C. (131xc2x0 F.), as measured by the Penske Martin Close Cup Method (ASTM D93), more preferably higher than about 60xc2x0 C. (140xc2x0 F.).\nThe inventors herein have found that a solvent having a vapor pressure less than about 5 mm Hg at 20xc2x0 C. must be used to avoid unacceptable latent losses in an insecticidal formulation over an extended period of delivery. Hydrocarbon solvents commonly used in analogous environments, such as Isopar C, Isopar G, Isopar L, Isopar E, and the like, have all been found to lead to unacceptable latent losses.\nDipropylene glycol dimethyl ether (DMM) has been found to meet the constraints of low vapor pressure (0.55 mm Hg at 20xc2x0 C.), low viscosity (1 cp at 25xc2x0 C.), and high water solubility as an indicator of low toxicity (35 g/100 ml water). Alone among the glycol ether solvents having these characteristics, DMM has a relatively high flash point (149xc2x0 F. (65xc2x0 C.)). Thus, in the preferred embodiments, DMM is the main solvent, used with one or more co-solvents.\nIn preferred embodiments, the insecticidal composition includes a co-solvent. The function of co-solvents is to improve the efficacy of the final formulation while still enabling the formulation to meet the constraints of high flash point, low volatility, low viscosity and low toxicity (as evidenced by high solubility in water). Accordingly the cosolvents themselves may have lower flash points, higher volatility, higher viscosity and lower solubility than the main solvent and still be acceptable for use in the solvent system according to the invention. Isopar L solvent, a hydrocarbon with very little water solubility and having an R65 hazard designation, may be used as a co-solvent in amounts less than about 10 percent by weight. Propylene glycol methyl ether acetate, which has a flash point of 114xc2x0 F. and also carries an R36 hazard designation may be used as a co-solvent in amounts less than about 20 percent by weight with respect to the final formulation. Propylene glycol n-propyl ether (PnP), which has a flash point of 119xc2x0 F. may be used as a co-solvent in amounts up to about 20 percent by weight. A preferred co-solvent is the cyclic amide N-methylpyrrolidone (NMP) used in amounts less than about 20 percent by weight.\nThe insecticides usable in the composition according to the invention are not particularly limited and include, without limitation, pyrethroid insecticides such as pyrethrin, allethrin, tetramethrin, resmethrin, etc.; organic phosphorus insecticides such as fenitrothion, ciafos, diazinon, etc.; carbamate insecticides such as carbaryl. The synthetic pyrethroids transfluthrin ((IR-trans)-2,3,5,6-tetraflurophenyl) methyl 3-(2,2-dichloroethyenyl)-2,3-dimethylcyclopropane carboxylate, a synthetic pyrethroid), and prallethrin ((RS)-2-Methyl-4-Oxo-3-(2-Propynyl) Cyclopent-2-enyl (1RS)-cis, transchrysanthemate), sometimes referred to as Etoc, are particularly preferred. Teflutlrin, pynamin, pynamin forte, and natural pyrethrum are all acceptable for use.\nIn some of the following examples, Pyrethrum 20 is used as an insecticide. Pyrethrum 20, available from Waxdale, is a mixture of pyrethrins in petroleum distillates and vegetable extracts, having 20 percent by weight active ingredients."}
{"text": "Raman spectroscopy is a powerful tool for quantitative analysis of the composition and concentration of a certain analyte within a sample. In the application of physiological detection, such an optical measurement is usually carried out in the near-infrared region. The Raman signal is often small, for instance the ratio of Raman signal strength to excitation signal strength is less than 10−10. Moreover the Raman signal is ultra-sensitive to the measurement conditions, such as laser fluctuation, optical bleaching, temperature variation, the changes in sample size and sample shape and optical alignment. Therefore internal standard method is usually applied during the analysis of the Raman signal. The basic principle of internal standard method is to measure a sample signal and a standard signal simultaneously (or nearly simultaneously) and their ratio, which is invariant from the measurement conditions, is used in the quantitation. Conventionally, a Grating-CCD (or Grating-photodiode-array) spectrometer is used as the detector in these systems. However, these systems are expensive and throughput is limited by the grating at required spectral resolution.\nRecent developments in Microelectromechanical systems (MEMS)-based spectrometer made it possible to use single detector to replace the CCD system in Raman signal detection system which significantly reduces the system cost. However, the throughput of such single-detector system is still limited by the use of grating. In order to achieve high throughput system, different designs have been disclosed. One of such design uses a sweeping light source, for instance a tunable laser, to obtain the Raman spectrum. Another approach uses a tunable filter, for instance an acousto-optical tunable filter, to obtain the Raman spectrum. Nonetheless, the costs of these new designs are not acceptable in constructing a home-used device for monitoring physiological parameters."}
{"text": "This invention relates generally to dynamoelectric machines, and in particular to a cover which is removably mounted to an endshield of a motor and which simultaneously seals against the endshield and provides unhindered access to fasteners for connecting the machine to an adjacent part, such as a pump.\nDynamoelectric machines such as electric motors have widespread use in converting energy between electrical current and mechanical motion. A substantial concern is that motors are frequently exposed to contaminants such as moisture and dust which damage or degrade the motors. Electrical component parts used in control and operation of the motors are particularly vulnerable to damage from water, which can produce short-circuits and safety hazards. Consequently, electronic components are mounted at locations where they may be protected from environmental contamination. For example, a motor for an above-the-ground swimming pool filter pump is typically placed outdoors, adjacent to the pool, where it is exposed to both rain and splashes. The motor has an electric power cord with wires that connect at a terminal board holding electronic components. The terminal board is mounted inside the motor behind a sealed cover to prevent moisture from reaching the terminal board or its components, thereby lowering the risk of contamination. Typically, the cover includes a steel plate and gasket which are mounted to an endshield of the motor using two or more small threaded fasteners. Unfortunately, removal and replacement of that cover increases the time required to access the terminal board for assembly or maintenance, thereby degrading efficiency.\nSome motors have a dome-type cover which encloses a portion of the endshield, but these covers have the drawback of requiring additional installation time during assembly. Either they fail to provide an effective seal so that a second, steel plate cover is still required, or alternatively they obstruct access to bolt fasteners used in connecting the motor to an adjacent part, such as a pump or support. In the latter instance, the dome-type cover must be removed for manipulating the bolt fasteners even though there is no need to access the terminal board or electronic components. That exposes the electronics to potential contamination during the time when the cover is removed. More importantly, a person may inadvertently damage electronic components when removing or re-installing the cover, or may inadvertently re-install the cover at an unsealing, mis-aligned orientation, causing subsequent leaks and damage.\nAmong the several objects and features of the present invention may be noted the provision of a dynamoelectric machine which is connectable to an adjacent part while being simultaneously protected from moisture or contaminants; the provision of such a machine which reduces the time required for connection to the adjacent part; the provision of such a machine which has an effective seal; the provision of such a machine which provides unhindered access to fasteners; the provision of such a machine which reduces risk of damage to electronic components; and the provision of such a machine which is economical.\nIn general, a dynamoelectric machine of the present invention is connectable to an adjacent part while being simultaneously protected from moisture or contaminants. The machine comprises a casing having a hollow interior defining a housing and a longitudinal axis. An endshield is positioned generally at an end of the casing and mounts at least one electrical component used in operation of the machine. The endshield has a wall projecting in a longitudinal direction and further has at least two openings for receiving fasteners to connect the machine to the adjacent part. The at least two openings define respective fastening locations on the endshield. A cover is removably mounted on the endshield and defines an enclosure to protect the at least one electrical component. The cover is configured for an overlapping fit with at least a portion of the wall of the endshield to prevent passage of moisture or contaminants into the enclosure. The cover is further configured to avoid covering the at least two openings when mounted on the endshield such that the fastener locations lie outside of the enclosure and are accessible for connecting the machine to the adjacent part while the cover remains mounted on the endshield.\nIn another aspect, a method according to the present invention protects a dynamoelectric machine from moisture or contaminants while connecting the machine to an adjacent part. The machine has a casing, two opposite endshields, and at least two fasteners connecting the endshields and extending beyond the casing for connection to the adjacent part. The method comprises the steps of mounting a cover on a portion of one of the endshields with a fit which overlaps a sidewall of the endshield to prevent passage of moisture or contaminants. The cover is configured to avoid covering the fasteners. The fasteners are manipulated to connect the machine to the adjacent part while the cover remains mounted on the endshield.\nOther objects and features of the present invention will be in part apparent and in part pointed out hereinafter."}
{"text": "1,1,2-Trichloro-1,2,2-trifluoroethane (known in the trade under the name F113) has been widely used in industry for cleaning and degreasing highly-varied solid surfaces (metal components, glasses, plastics, composites) for which the absence or at least the lowest possible residual content of impurities, in particular of organic nature, is required. F113 was particularly well suited to this use because of its nonaggressive nature with regard to the materials used. This product was used in particular in the field of the manufacture of printed circuits, for removing the residues of the substances used to improve the quality of the soldered joints (denoted by the term solder flux). This removal operation is denoted in the trade by the term \"defluxing\".\nMention may also be made of the applications of F113 in the degreasing of heavy metal components and in the cleaning of mechanical components of high quality and of great accuracy, such as, for example, gyroscopes and military, aerospace or medical equipment. In its various applications, F113 is generally used in combination with other organic solvents (for example methanol), in order to improve its cleaning power. It is then preferable to use azeotropic or near-azeotropic mixtures. The term \"near-azeotropic mixture\" is understood to mean, within the sense of the present invention, a mixture of generally miscible chemical compounds which, under certain specific conditions of proportions, temperature and pressure, boils at a substantially constant temperature while retaining substantially the same composition. When it is heated to reflux, such a near-azeotropic mixture is in equilibrium with a vapour phase, the composition of which is substantially the same as that of the liquid phase. Such azeotropic or near-azeotropic behaviour is desirable in ensuring satisfactory operation of the devices in which the abovementioned cleaning operations are carried out and in particular in ensuring the recycling by distillation of the cleaning fluid.\nF113 is also used in fields, in particular in optics, where it is required to have available surfaces which are devoid of water, that is to say surfaces where water is only present in the form of traces undetectable by the measurement method (Karl Fischer method). F113 is, for this purpose, employed in drying (or dewetting) operations on the said surfaces, in combination with hydrophobic surface-active agents.\nHowever, the use of compositions based on F113 is now forbidden as F113 is one of the chlorofluorocarbons (CFCs) suspected of attacking or damaging the stratospheric ozone.\nIn these various applications, F113 can be replaced by 1,1-dichloro-1-fluoroethane (known under the name F141b), but the use of this substitute is already controlled because, although low, it still has a destructive effect with regard to ozone.\nApplication EP 0,512,885 discloses a composition, comprising from 93 to 99% by weight of 1,1,1,3,3-pentafluorobutane and from 1 to 7% of methanol, which can be used as substitute for F113. 1,1,1,3,3-pentafluorobutane, also known in the trade under the name F365mfc, has no destructive effect with regard to ozone.\nApplication EP 0,856,578 discloses a composition, comprising from 10 to 90% by weight of 1,1,1,2,3,4,4,5,5,5-decafluoropentane, from 10 to 90% of dichloromethane and from 0 to 10% of methanol, which can also be used as substitute for F113. 1,1,1,2,3,4,4,5,5,5-Decafluoropentane, known in the trade under the name 43-10mee, also has no destructive effect with regard to ozone."}
{"text": "Mainframe manufacturers and the telecommunications industry are primary users of backplanes on which are mounted a substantial number of daughter cards carrying electronic components thereon. The backplane provides conductive traces or circuits to electrically interconnect the components on the several cards and to provide access to other backplanes and outside electric gear. The predominate method of mounting the cards on the backplane and electrically interconnecting the circuits (and hence the components) is by means of card edge connectors. Such connectors include a card receiving slot, i.e. a card slot, with conductive contact members having spring arms or cantilever beams positioned along one or both sides of the slot to engage traces on the card inserted therein and further having depending leads electrically engaging the backplane. In some connectors, the arms or beams are preloaded into the slot and the cards are frictionally inserted thereinto. Provided the number of contact members are not too great, the force required to insert the card is acceptable. However, cards having a large number of traces require connectors wherein the large number of contact members prohibit frictional insertion. In these cases, connectors having camming mechanisms which either cam the arms or beams out of the slot prior to inserting the card, e.g. as disclosed in U.S. Pat. No. 3,793,609, or into the slot after insertion of the card, e.g. as disclosed in U.S. Pat. No. 4,586,772, are used.\nIn each type connector, the method of obtaining the normal force, i.e. the force exerted on the conductive pads or traces on the card by the arms or beams, predetermines the type of cantilever beam or spring arm to be used therein. For example, in the first type, the spring arms must be resilient enough so that they can be moved out of the slot with acceptable levels of force applied to the camming mechanism and still have sufficient inherent spring force to bear against the card and obtain good electrical contact therewith. In the second type of ZIF connector, the spring arms or beams must be made of stiffer material to stand up under the biasing forces continually exerted against them by the camming mechanism. In this type connector, a higher normal force may be obtained which is required in some operational uses. However, a greater force is required to actuate the camming mechanism which could lead to breakage if the components are not strong enough to withstand the forces.\nIt is now proposed to provide a ZIF card edge connector wherein the spring arms on the contact members are hinged and have a large moment arm which provides a mechanical advantage so that lower forces will actuate the camming mechanism and a high normal force will still be obtained."}
{"text": "Photovoltaic cells are commonly used to transfer energy in the form of light into energy in the form of electricity. A typical photovoltaic cell includes a photoactive material disposed between two electrodes. Generally, light passes through one or both of the electrodes to interact with the photoactive material, thereby generating charge carriers (i.e., electrons and holes). As a result, the ability of the photoactive material to absorb light and general charge carriers can limit the overall efficiency of a photovoltaic cell."}
{"text": "In recent years, with rapid spread of information-related equipment such as personal computers, video cameras, and mobile phones, batteries used as a power source therefor have been actively developed. Further, in the automobile industry, development of batteries having high output and high capacity for electric vehicles or hybrid vehicles has been promoted. Among various batteries, a lithium battery has an advantage of a high energy density.\nIn this regard, as a solvent of a liquid electrolyte for a lithium battery, it is known to use a glyme that is one of oligoethers. For example, Patent Literature 1 discloses a Li-based secondary battery which includes a cathode containing a specific electrode active material, a liquid electrolyte in which at least a part of a glyme and a Li salt forms a complex, and an anode. This technique is intended to improve rate characteristics and a power density of a battery. In addition, Patent Literature 2 discloses a liquid electrolyte which contains only a glyme as a solvent of the liquid electrolyte and an alkali metal salt as a solute of the liquid electrolyte, in which at least a part of the glyme and the alkali metal salt forms a complex. This technique is intended to provide a high-voltage liquid electrolyte having high thermal stability and being usable in general in-vehicle secondary batteries."}
{"text": "Industrial coatings are surface protective coatings (paint coatings) applied to substrates and typically heat cured to form continuous films for decorative purposes as well as to protect the substrate. A protective coating ordinarily comprises an organic polymeric binder, pigments, and various paint additives, where the polymeric binder acts as a fluid vehicle for the pigments and imparts rheological properties to the fluid paint coating. Upon curing, the polymeric binder hardens and functions as a binder for the pigments and provides adhesion of the dried paint film to the substrate. The pigments may be organic or inorganic and functionally contribute to opacity and color in addition to durability and hardness, although some paint coatings contain little or no opacifying pigments and are described as clear coatings. The manufacture of paint coatings involves the preparation of a polymeric binder, mixing of component materials, grinding of pigments in the polymeric binder, and thinning to commercial standards.\nEpoxy resins are particularly desirable for use in protective surface coating materials as a vehicle or polymeric binder for pigments, fillers, and other additives where the epoxy resins advantageously provide toughness, flexibility, adhesion, and chemical resistance. Water-dispersed coating compositions containing epoxy resins are highly desirable for can coating compositions and particularly useful for interior surfaces. Comings for soft drink and beer cans, for instance, are critical due to taste sensitivity wherein such can coatings must not alter the product taste of canned beverages. Taste problems can occur in a variety of ways such as by leaching of coating components into the beverage, or by adsorption of flavor by the coating, or sometimes by chemical reaction, or by some combination thereof.\nContainer coating technology frequently utilizes an epoxy resin which has been grafted with acrylic monomers, styrene, and methacrylic acid. This grafted epoxy resin is prepared in solvent, usually butyl cellosolve, and n-butanol, to maintain low processing viscosities and then reduced with water by a direct or inverse let down procedure. Although cured film properties are highly desirable, such coatings suffer from the fact that sizeable amounts of solvent are required to obtain good performance. High molecular weight epoxy resins typically require 50% to 90% solvent (based on total solids plus organic solvent) before reducing with amine and water.\nEpoxy based can coatings comprising a carbon grafted acrylic chain produced in the presence of an extender resin are disclosed in U.S. Pat. No. 4,399,241 and U.S. Pat. No. 4,482,671 while U.S. Pat. No. 4,595,716, and U.S. Pat. No. 5,157,078 teach a carbon grafting process involving solvent polymerization at moderate temperatures with high levels of peroxide initiator to produce a carbon-graft polymer. The high solvent levels, however, invariably carry over to the aqueous dispersion when the resulting polymers are dispersed into water to produce a VOC (volatile organic compounds) level considerably above 2 and typically between 3 and 4 pounds volatile organic compounds per gallon of resin solids.\nCommonly assigned U.S. Pat. No. 5,290,828 discloses an acrylic grafted epoxy polyester terpolymer produced by in-situ copolymerization of ethylenic monomers with low molecular weight epoxy and unsaturated polyester resins where carboxyl monomers esterify epoxy groups while monomer double bonds coreact with polyester double bonds to form the terpolymer. Commonly assigned Ser. No. 222,029 filed Apr. 4, 1994 discloses an acrylic grafted epoxy-ester produced by first esterifying epoxy resin with a carboxyl functional unsaturated polyester to form an unsaturated epoxy-ester. The unsaturated epoxy-ester is dispersed into water and followed by in-situ copolymerization of ethylenic monomers in the aqueous dispersion, where the copolymerized monomers partially graft to the preformed unsaturated epoxy-ester.\nCoating compositions based on microgels are shown in U.S. Pat. No. 4,897,434 where epoxy esters are preformed, then dispersed into water, and thereafter further crosslink available epoxy and carboxyl groups on the preformed epoxy-ester.\nIt now has been found that a preformed, carboxyl functional, low molecular weight acrylic addition copolymer of copolymerized ethylenic monomers, including carboxyl functional ethylenic monomer, can be utilized advantageously as an aqueous dispersing or suspending agent for dispersing an epoxy resin into water to produce a very low VOC coating having a VOC (volatile organic compounds) less than about one pound per gallon of resin solids. The resulting water dispersed epoxy resin provides an emulsion polymerization medium, where the addition copolymer dispersing agent is esterified with the epoxy functional resin to form crosslinked microgel particles stably dispersed in water. The acid functional, low molecular weight addition copolymer, preferably an acrylic copolymer, (preferably 2,000 to 20,000 Mn) can be dispersed in ammonia water to form small particle size dispersions (ca 50 nm), while subsequent crosslinking with liquid epoxy resin converts these dispersions to microgels. The water dispersed, crosslinked small particle size provide good storage stability, rheology control, along with exceptional film properties including excellent water resistance, low temperature cure, excellent flexibility, and good resistance to odor adsorption. This invention avoids the use of acid functional polyesters, which can slowly hydrolyze causing gelation while rheology of very low viscosity products can only be controlled by adding rheological control agents. In this invention, the acrylic copolymer dispersant can be synthesized at a very low solvent level (&lt;1 lb/gal VOC), dispersed into ammonia water, and subsequently crosslinked with liquid epoxy resin to form stable microgel particles. Crosslinking of the addition copolymer with diepoxide while stably dispersed in water produces very small size crosslinked microgel particles, a physical property particularly useful for producing tough but resilient coating films applied to a substrate. Further advantages rely on an acid functional acrylic to produce small particle size dispersions (ca. 50 nm) with low applied shear which will not hydrolyze and provide excellent storage stability. Viscosity can be easily controlled by adjusting acrylic Mn, carboxylic acid concentration, comonomer type, and degree of reaction with epoxy resin. These and other advantages of this invention will become more apparent by referring to the detailed description and illustrative examples."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display unit and a method for its production.\nIn recent years, the display areas of liquid crystal display units of such devices as computer terminals, personal computers, word processors, TVs and the like have become increasingly larger, and their uses are expected to become further expanded in the future, accompanied by greater demand to realize liquid crystal display units with visual angle characteristics equaling those of CRTs (cathode ray tubes).\n2. Description of the Related Art\nSome techniques have already been proposed for improving viewing angle characteristics, including a PDN type (Polymer Dispersed Liquid Crystal Display with Crossed Nicols; see Japanese Unexamined Patent Publication No. 4-212928), an amorphous TN type (Amorphous Twisted Nematic Liquid Crystal Display; see SID '94 Digest, p. 915-918) and a PTN type (Polymer Twisted Nematic Liquid Crystal Display; see Japanese Unexamined Patent Publication No. 5-27242).\nFIGS. 7A and 7B are illustrations of PDN-type liquid crystal display units of the prior art, FIG. 7A showing a state without application of voltage and FIG. 7B showing a state with application of voltage.\nIn these drawings, 21 is a transparent base plate, 22 is a transparent conductive film, 23 is a transparent base plate, 24 is a transparent conductive film, 25 is a polymer, 26 is a liquid crystal molecule, 27 is a liquid crystal droplet, 28 is a liquid crystal layer, 29 and 30 are polarizer, 31 is a round cushion-shaped film with optical anisotropy, 32 is an alternating current power source and 33 is a switch.\nIn this conventional PDN-type liquid crystal display unit, the opposing transparent base plate 21 with the transparent conductive film 22 and transparent base plate 23 with the transparent conductive film 24 form a structure sandwiching the liquid crystal layer 28 containing the polymer 25 and liquid crystal droplets 27. Also, the pair of polarizers 29 and 30 whose directions of polarization are orthogonal are laid over the outside of the transparent base plates 21 and 23.\nA round cushion-shaped film 31 with optical anisotropy, such as a polycarbonate polyaxially stretched film, is laid between the transparent base plate 23 and the polarizer 30 on the display side.\nThe alternating current power source 32 and the switch 33 are connected in series between the transparent conductive film 22 and the transparent conductive film 24.\nIn this liquid crystal display unit, when no voltage is applied from the alternating current power source 32 between the transparent conductive films 22 and 24, the liquid crystal molecules 26 in each of the liquid crystal droplets 27 are controlled so as to line up at the interface with the polymer 25, but since the orientation of the liquid crystal molecules 26 on the liquid crystal droplets 27 is unordered, incident light 34 from the back side is scattered by the liquid crystal molecules 26 and passes through the polarizing plate 30 appearing milky white, and bright (see FIG. 7A).\nWhen a voltage is applied from the alternating current power source 32 between the transparent conductive films 22 and 24, the liquid crystal molecules 26 in the liquid crystal droplets 27 align in the direction of the electric field, and as a result light scattering is reduced and incident light 34 from the back side may pass through the liquid crystal layer 28 unhindered.\nHowever, since the polarizers 29 and 30 are arranged in crossed Nicols, the light cannot pass through the polarizer 30, giving a dark appearance.\nWhen the liquid crystal panel of this liquid crystal display unit is viewed at an oblique inclination angle, the liquid crystal molecules 26 and the incident light 34 both are viewed from an oblique inclined angle and the optical effect of the liquid crystal layer 28 appears to some degree causing light leakage, and therefore the round cushion-shaped film 31 with optical anisotropy is inserted to block the resulting leaking light which passes through the polarizing plate 30.\nGenerally speaking, when the refractive index of a liquid crystal molecule is expressed in three-dimensional terms it may be compared to a rugby ball-shaped index ellipsoid, and the refractive index when viewed at an oblique inclination angle corresponds to the ellipse section resulting from cutting the rugby ball-shaped index ellipsoid at an oblique inclination angle.\nOn the other hand, when the refractive index of a molecule of the round cushion-shaped film with optical anisotropy is expressed in three-dimensional terms, it exhibits the characteristics of a round cushion-shaped index ellipsoid, and the refractive index thereof when viewed at an oblique inclination angle corresponds to the ellipse section resulting from cutting the round cushion at an oblique inclination angle.\nAlso, the ellipse section resulting from cutting the rugby ball-shaped index ellipsoid of the liquid crystal molecule at an oblique inclination angle and the ellipse section resulting from cutting the index ellipsoid of the round cushion-shaped film with optical anisotropy at an oblique inclination angle have their long axes and short axes rotated 90.degree. to each other.\nConsequently, when the ellipse section resulting from cutting the rugby ball-shaped index ellipsoid of the liquid crystal molecule at an oblique inclination angle and the ellipse section resulting from cutting the index ellipsoid of the round cushion-shaped film with optical anisotropy at an oblique inclination angle are superimposed a circle results, and the optical effect is canceled out.\nTherefore, as mentioned above, the superimposition of the round cushion-shaped film 31 with optical anisotropy cancels out the effect of the inclined incident light, and therefore the black display produced when a voltage is applied between the electrodes has a consistently dark appearance even when viewed from an oblique inclination angle, giving a constant display when viewed from any direction, and producing the effect of a wider visual angle.\nFIGS. 8-A-1 to 8-A-4 and 8-B-1 to 8-B-4 are illustrations of amorphous TN-type liquid crystal display units of the prior art, FIGS. 8-A-1 to 8-A-4 showing the prior art and FIGS. 8-B-1 to 8-B-4 showing amorphous TN-type liquid crystal display units.\nIn these drawings, 41 is a first transparent conductive base plate, 42 is a second transparent conductive base plate, 43 and 44 are alignment films, 45 is a liquid crystal molecule, 46 is a liquid crystal layer and 47 and 48 are non-rubbed polymer films.\nA prior art TN (Twisted-nematic) liquid crystal display unit with rubbed alignment films will now be explained with reference to FIGS. 8-A-1 to 8-A-4.\nIn this TN liquid crystal display unit of the prior art, the liquid crystal layer 46 comprised of liquid crystal molecules 45 is sandwiched between the first transparent conductive base plate 41 with an alignment film 43 which is rubbed in the direction perpendicular to the plane of the paper, and the second transparent conductive base plate 42 with an alignment film 44 which is rubbed left to right on the plane of the paper.\nIn this TN liquid crystal display unit, even at clearing temperatures at which the orientation of the liquid crystal molecules 45 is random (isotropic), the liquid crystal molecules 45 contacting with the alignment films 43 and 44 of the first transparent conductive base plate 41 and the second transparent conductive base plate 42 which are rubbed in orthogonal directions are aligned in these respective directions of alignment by the restraint of the alignment films 43 and 44 (see FIG. 8-A-1).\nWhen the temperature is lowered below the clearing temperature, the effect of the liquid crystal molecules 45 contacting with the upper and lower alignment films 43 and 44 reaches inside the liquid crystal layer 46, and a column of liquid crystal molecules 45 begins to grow (see FIG. 8-A-2).\nWhen the temperature is further lowered, the effect of the liquid crystal molecules 45 contacting with the upper and lower alignment films 43 and 44 reaches further inside the liquid crystal layer 46, and the column of liquid crystal molecules 45 grows larger (see FIG. 8-A-3).\nWhen the temperature is even further lowered, the effect of the liquid crystal molecules 45 contacting with the upper and lower alignment films 43 and 44 reaches even further inside the liquid crystal layer 46, and the column of liquid crystal molecules 45 is completely formed between the upper and lower alignment films 43 and 44. Since the twist origins of the liquid crystal molecules 45 are the same in state, the viewing angle characteristics are narrower (see FIG. 8-A-4).\nAn amorphous TN-type liquid crystal display unit will now be explained with reference to FIGS. 8-B-1 to 8-B-4.\nIn this amorphous TN-type liquid crystal display unit, the liquid crystal layer 46 comprising liquid crystal molecules 45 is sandwiched between a first transparent conductive base plate 41 and a second transparent base plate 42 with non-rubbed polymer (polyimide PI) films 47 and 48.\nIn this amorphous TN-type liquid crystal display unit, at clearing temperatures at which the orientation of the liquid crystal molecules 45 is random, the surfaces of the non-rubbed polymer films 47 and 48 in contact with the liquid crystal molecules 45 are rough, and therefore the liquid crystal molecules 45 are not aligned, being oriented in random directions (see FIG. 8-B-1).\nWhen the temperature is further lowered, a column of aligned liquid crystal molecules 45 does not grow from the surface of the liquid crystal layer 46 in contact with the non-rubbed polymer films 47 and 48, but rather grows from the interior (bulk region) (see FIG. 8-B-2).\nWhen the temperature is even further lowered, the alignment of the liquid crystal molecules 45 occurring in the bulk region is promoted, causing the column of liquid crystal molecules 45 to grow larger (see FIG. 8-B-3).\nWhen the temperature is still further lowered and a temperature of 30.degree. C. below the clearing temperature is maintained, the alignment of the liquid crystal molecules occurring in the bulk region is further promoted, and the column of liquid crystal molecules 45 in each region is completed between the non-rubbed polymer films 47 and 48, with the liquid crystal molecules 45 becoming rigidly fixed in that state (memory effect).\nIn this state, the twist origins of the liquid crystal molecules 45 in contact with the surfaces of the non-rubbed polymer films 47 and 48 are random, and therefore the viewing angle characteristics of each of the columns of liquid crystal molecules 45 are averaged to provide improved viewing angle characteristics and to obtain a gray scale (see FIG. 8-B-4).\nFIG. 9 is an illustration of a PTN-type liquid crystal display unit of the prior art.\nIn this drawing, 51 and 52 are transparent conductive base plates, 53 denotes liquid crystal molecules, and 54 denotes a material which disturbs the twisted arrangement of the liquid crystal molecules.\nIn this conventional PTN-type liquid crystal display unit, liquid crystals such as cholesteric liquid crystals or chiral nematic liquid crystals are injected between the pair of transparent conductive base plates 51 and 52 with transparent electrodes and alignment films on their opposing inner sides, and there is also injected a material 54 to disturb the twisted arrangement of the liquid crystal molecules, for example an organic polymer such as polymethyl methacrylate, polystyrene melamine resin, urea resin, phenol resin or polydiisopropyl fumarate, or a photopolymerized polymer such as polyurethane acrylate, polyester acrylate, epoxy acrylate or polyether acrylate, in order to change the arrangement of the liquid crystal molecules 53.\nAs a result, the twisted orientations of the liquid crystal molecules along the lines A, B and C vary slightly depending on the position of the liquid crystal layer, even with the same viewing angle direction.\nSince the twisted orientations of the liquid crystal molecules 53 vary slightly along the lines A, B and C in this drawing, application of a voltage causes the liquid crystal molecules 53 to be aligned differently, and the orientations of the liquid crystal molecules 53 are averaged with respect to the viewing angle direction, thus resulting in averaging of the apparent transmittance/applied voltage characteristics, improved viewing angle characteristics, and suppression of the abnormal display inversion phenomenon.\nIn the PDN-type liquid crystal display unit described above it has not been possible to impart optical activity to the liquid crystal molecules, and consequently the display brightness of displayed white has been only 1/3 that of TN-types.\nIn addition, the sizes of the domains lack uniformity in amorphous TN-type liquid crystal display units and are large, and when the display surface is viewed from an oblique inclination angle a patchy black and white pattern is produced.\nWhen the display surface is viewed from an oblique inclination angle, the black portions and white portions are of an extremely small identical size and there are few problems so long as they are dispersed, but since the domains are large, they appear as patches.\nThis condition is the same in the case of PTN-type liquid crystal display units, which have had the problem of patchiness when viewed from an oblique inclination angle.\nIt is an object of the present invention to provide a liquid crystal display unit with satisfactory viewing angle characteristics and which eliminates such patches when viewed from an oblique inclination angle."}
{"text": "Transactions are familiar to us in everyday life when we exchange money for goods and services, such as buying a train ticket or paying for medicine. Such transactions involve a short conversation (for example, requesting availability and cost), and then making the payment. The processing of any one of these items is a business transaction of the kind that is handled by a transaction management system.\nA transaction in the business sense can be viewed as an activity that must be completed in its entirety with a mutually agreed-upon outcome. It usually involves operations on some shared resources and results in overall change of state affecting some or all of those resources. When an activity or a transaction has been started and the mutually agreed outcome cannot be achieved, all parties involved in a transaction should revert to the state they were in before its initiation. In other words, all operations should be undone as if they had never taken place.\nThere are many examples of business transactions. A common one involves transfer of money between bank accounts. In this scenario, a business transaction would be a two-step process involving subtraction (debit) from one account and addition (credit) to another account. Both operations are part of the same transaction and both must succeed in order to complete the transaction. If one of these operations fails, the account balances must be restored to their original states.\nTypically such transactions consist of many computing and data access tasks to be executed in one or more machines; the tasks may include handling the user interface, data retrieval and modification, and communications. In the example above, the money transfer operation is a transaction composed of debiting one account and crediting another.\nIn the context of business software, we can express the above more precisely. A transaction (sometimes also referred to as a ‘Unit-Of-Work’ or ‘UOW’) is a set of related operations that must be completed together. All their recoverable actions must either complete successfully or not at all. This property of a transaction is referred to as ‘atomicity’.\nIn the simplest case, a transaction will access resources connected to a single computer processor. Such a transaction is referred to as a ‘local transaction’. More often, however, a transaction will access resources which are located in several different computer processors, or logical partitions, each maintaining its own transaction log. Such a transaction is referred to as a ‘distributed transaction’.\nWhen a distributed transaction ends, the atomicity property of transactions requires that either all of the processors involved commit the transaction or all of them abort the transaction. To achieve this one of the processors takes on the role of coordinator to ensure the same outcome at all of the parties to the transaction, using a ‘coordination protocol’ that is commonly understood and followed by all the parties involved. The two-phase commit protocol has been widely adopted as the protocol of choice in the distributed transaction management environment. This protocol guarantees that the work is either successfully completed by all its participants or not performed at all, with any data modifications being either committed together (completed) or rolled back (backed out) together on failure.\nAnother property of a transaction is its durability. This means that once a user has been notified of success, a transaction's outcome must persist, and not be undone, even when there is a system failure. A recovery manager is used to ensure that a system's objects are durable and that the effects of transactions are atomic even when the system crashes. The recovery manager saves information required for recovery purposes. This recovery can be for the dynamic backout of a transaction, perhaps as a result of a failure after a task updated a recoverable temporary storage queue. Additionally the recovery data can be used for restoring a transaction processing system to a committed state, for example when the system is restarted after system failure. Typically, the recovery file comprises at least one log containing the history of the transactions performed by a transaction processing system. In the event of a system failure, the recovery file can be played back to return the system to its state right before the failure, and the transaction log(s) used to check for and undo transactions that were not properly completed before failure.\nAlso, in the event of a transaction failure, the transaction log can be used to reverse updates that have been carried out by that transaction, by working backwards from the last change before the failure, hence the name ‘dynamic transaction backout’. The backout occurs before any locks on any affected resources are released, which safeguards other tasks from the possibility of using corrupted data, because modified data is not released for use by them (“committed”) until the current task has finished with it. In case the log needs to be replayed later in a system restart, an entry is first made in the log to indicate that that transaction is being backed out.\nExamples of systems which carry out such transaction logging include transaction systems such as IBM® CICS® Transaction Server or IBM WebSphere® Application Server, as well as database systems such as IBM DB2® or IBM IMS™ (IBM, CICS, WebSphere, DB2 and IMS are trade marks of International Business Machines Corporation).\nThe scheduling and dispatching of transaction processing tasks is usually handled by a dispatcher, which controls the order in which tasks are handled by allocating processor time among them, based on a dispatching priority value associated with each task. The log typically records the information in the order that the activity occurs. Without some management, this would consume an ever increasing amount of resource. So it must be reorganised on a regular basis so as to reduce its size by the recovery manager carrying out a process called ‘keypointing’. Keypointing comprises writing current committed values of a system's objects to a new portion of the recovery file, together with transaction status entries and intentions lists for transactions that have not been fully resolved. An intentions list for a transaction contains a list of the references and the values of all the objects/resources that are altered by that transaction, as well as information related to the two-phase commit protocol. Once a keypoint has been taken, i.e. information stored through a keypointing procedure, recovery information for irrevocably committed (or backed out) transactions can usually be discarded from the log, sometimes called ‘trimming’ the log. This reduces the file size of the log as well as the number of transactions to be dealt with during recovery.\nThe rate that old log records are deleted by trimming should ideally match the rate at which new log data is being generated as new work enters the transaction system and makes its own recoverable changes to resources. In a well tuned on-line transaction processing (OLTP) environment, the volume of log data held on the log should remain reasonably constant over time, with an underlying periodic rise and fall in the number of log records held on the log as both the new work and the housekeeping work are performed.\nThis mechanism of log data deletion usually solves the problem of logs continually growing in size over time. However, a long-running UOW can prevent this mechanism from working, as until such a UOW has completed, and its log data is no longer required, its data cannot be deleted from the log and all data logged since the first log entry of that UOW must be maintained on the log. Hence, the system will not be able to trim the log after keypointing, and the log will grow and grow in size. A long-running unit of work (UOW) of this type may be a valid UOW that has performed a large number of recoverable changes, and so generated a large volume of log data. Alternatively, it may be for a so-called ‘rogue’ program that has entered an invalid loop while performing continual recoverable changes, thereby generating an excessive number of log records, or which is hanging in an uncommitted state.\nEventually, when a critical threshold of the logging subsystem, the operating system, or the available hardware is exceeded, a request to write to the log will fail for an “out of space” type of condition. When such an error is returned to the transaction system, it typically results in a serious failure. The system can no longer log any recoverable changes and so protect them from failures that require them to be backed out. Recovery processing (and hence data integrity) can no longer be guaranteed.\nCurrent mechanisms intended to prevent such a scenario include the emission of human-readable messages when a certain number of keypoints have been taken without the system having been able to trim the log. Such messages typically indicate the UOW which is preventing the log from being trimmed. In response to such a message, or to a series of such messages, a system operator may choose to take manual action to terminate that ‘rogue’ UOW.\nThe present invention seeks to address the above problems and provide an improved mechanism which may reduce the likelihood of encountering a log-full condition."}
{"text": "The invention relates to a system for inspecting and/or processing a sample on a stage with a probe or imaging apparatus.\nHigh resolution lithography and micro-machining tools (e.g. E-beam lithography system, DUV-stepper), inspection and process-control equipment (e.g. CD-SEM, defect location and defect review tools) and analytical and test equipment (SEM, FIB, AFM, high resolution tools using UV, DUV, and EUV light sources) require stages for carrying and moving the sample to be inspected and/or processed.\nWith increasing resolution requirements of those systems, the demands on these stages with respect to stability and position accuracy are steadily increasing. Accuracy and stability requirements are below 100 nanometer and will probably go down below 1 nanometer. The mechanical design and mechanics of the stage no longer can fulfill the challenging demand alone. Therefore, it is necessary to apply position control means to support the stage position accuracy and stage stability.\nStage position control is performed mainly by two tools, by a linear optical encoder and by a laser-interferometer. A linear encoder consists of optical grids or gratings which are connected to the stage and whose motion is read by suitable optical sensors, which are placed on the non-moveable part of the stage setup. The laser-interferometer control uses the interference between a reflected beam, which is generated by a mirror connected with the stage, and a reference beam. Since the wavelength of the laser beam is used as a measurement unit, laser interferometers offer a much better position accuracy (in the range of 10 nm) than optical encoders which have an accuracy in the range of 100 nm.\nLinear encoders have the drawback that their accuracy is limited. Laser-interferometers are much more accurate, but they are very complex and expensive. The accuracy is limited by changes in the refractive index of the optical light path and the quality of the mirrors. Furthermore, laser-interferometers are used preferably for x-y-coordinate measurements. Setups for a 3-dimensional coordinate measurement of stage position is very difficult to realize.\nFrom EP-A-0 559 402 and U.S. Pat. No. 3,936,716 stabilization systems are known for stabilizing a platform with gyroscopic means.\nIt is an object of the invention to provide a system for inspecting and/or processing a sample on a stage with a probe or imaging apparatus with a precise, simple and cost-effective position control means.\nIt is a further object of the invention to provide means which allows both the accurate control of the position of the stage relatively to the probe or imaging apparatus and the mechanical detection and stabilization of mechanical disturbances (e.g. vibrations, drift).\nThe position control means can be preferably used in probe systems (particle beam systems, atomic force microscopes, high resolution tools).\nThe gyroscopic means is adapted to measure movements in one, two or three dimensions.\nIn a preferred embodiment, there is at least one motor for moving the stage and/or the probe or imaging apparatus which is operatively connected with the position control means.\nThe system for inspecting and/or processing a sample according to a preferred embodiment comprises a probe or imaging apparatus. The probe and/or imaging apparatus according to the disclosed embodiment is defined by a scanning probe or a column with beam generating means and means to deflect the beam. In case of a column, the column is equipped with column gyroscopic means to detect movements of the optical column. Furthermore, the control means is operatively connected with the deflection means in order to compensate uncontrolled movements (vibrations, drift) of the stage relatively to the column.\nIn another system for inspecting and/or processing a sample, the probe or imaging apparatus is defined by an optical column having beam generating means, means to deflect the beam and image generating means wherein the control means is operatively connected with the image generating means in which the unintentional movements of stage relative to the column will be taken into account when processing the image."}
{"text": "The invention concerns a method and an apparatus which is especially advantageous for its performance in the thermal longitudinal parting of rectangular metal plates, in particular of cut-to-length continuous casting plate slabs while the metal plates rest on supports which run crosswise to the longitudinal direction of the parting. In the cross-parting of steel, cast as a continuous slab, or of NE-metals, the thermal separation on the basis of oxygen-flame cutting or plasma-flame cutting serves the preparation for the rolling which follows, whereby the cross-parting determines the length of, for instance, the coil weight of the tin sheeting to be rolled. Compared with that, the conditions, when parting metal slabs longitudinally, present themselves differently because in the latter case the desired width of the tin sheeting to be rolled out is determined. In recent years the aim was to increase the production rate in continuous casting through increasing the batch weights and through larger slab widths. Wider slabs do allow for a more economical production of casting products under maximal exploitation of the slab-casting apparatus. However, the wider slabs must be reduced to existing first-pass cross sections of the rolling mills, the width dimensions of which in turn go necessarily by the widths of the tin sheeting which are required in the subsequent procedure.\nIt is well-known (DE-AS No. 27 10 564) how, in order to fulfill this requirement, to move the cutting arc relative to the plate slab when parting crosswise, in which case, according to this suggestion vis-a-vis parting lengthwise the only difference is that the plate slab is at rest. The well-known longitudinal parting procedure requires, however, special steps for the protection of the supports which pass in crosswise direction to the longitudinal parting direction. Therefore, a diversified formation is required of the existing supports which go crosswise to the longitudinal parting direction. At first, individually driven rolls cause the transportation of the plate slabs. These rolls must be able to swing out in order not to be destroyed during the flame cutting. Further, crosswise-passing rigid plate slabs are provided which have to have so-called burn point openings for the path of the cutting arc. The burn point openings which are arranged in the path of the cutting arc result in a slit-like burn canal, the length of which corresponds to the length of the plate slabs, for instance 8 m and more, turns out to be very extended. A further disadvantage is that the slag and an increased amount of gas must be carried away from the long burn canal. Besides, the burning canal is firmly built-in and does not allow for any adjustment to changing the widths of the plate slabs to be parted longitudinally.\nThe invention is based on the task to simplify the known method, and based on that to create a better protected apparatus with less expenditure.\nThis task is solved in accordance with the invention by moving the metal slabs in longitudinal parting direction vis-a-vis the cutting arc which is motionless in the longitudinal parting direction and which is in a given case movable crosswise to the longitudinal parting direction only because of corrections in width conditional on deviations in width or through skewed position of the metal slabs, and by carrying off between two supports the slag collected at some point crosswise to the longitudinal parting direction between two supports. This proposal follows neither the principle of crosswise parting, in which case the continuously cast slab is moved by the pouring velocity which is followed by the cutting arc, nor does it follow the known principle where the continuous-cast plate slab is at rest. Consequently, it is advantageous that neither swing-away support rolls nor a relatively long burn-canal are necessary but that a burn-canal results naturally between the two existing support rolls. Accordingly, the burn-canal is short and the collected burn slag may be removed via a short passage. Furthermore, there is a resultant decrease in the volume of gas-exhaust and accordingly less power for the drive of the gas-exhaust. Another advantage of the invention lies in the fact that heavy crosswise-slidable slab centering machines are avoided, which are replaced by control arrangements which move the much lighter cutting arc installations. A further advantage lies in the possibility of changing the width-program without displacing the burn-point openings, i.e., to insert partings of 1/2, 1/3, 1/4 and the like for short runs.\nAn increased cutting performance may now be achieved after a further improvement of the invention, in that at least two metal slabs are moved relative to motionless cutting arcs which are at any given time spaced in the longitudinal parting direction. The doubling of the cutting capacity leads to an increased yield in cut goods and to an increased performance in comparison with cutting machines which run counter and thereby necessarily impede each other on encounter.\nThe performance of longitudinal parting is increased by an additional proposal in method when at a given time a metal slab is being cut only partially in the longitudinal direction, the cutting arc switched off afterwards, and this same metal slab is completely parted under a cutting arc which is pre-positioned in the longitudinal parting direction at a distance of more than the length of a slab, while the following metal slab is again parted, in longitudinal direction, only partially with the cutting arc again switched on.\nThe apparatus for the application of the invented method is designed in such a way that a cutting arc is provided between two rigidly arranged roller-bed rolls supporting the metal slabs along their full width. The cutting machine is motionless in the longitudinal parting direction and may be adjusted crosswise to the longitudinal parting direction. A pair of driving rolls are mounted adjacent the cutting machine. A conveyor duct for the resultant slag is provided, arranged crosswise to the longitudinal parting direction under the cutting machine, and a gas-exhaust pipe is connected to the conveyor duct. No swing-away rolls or additional supports are needed. The cutting machine, which is only movable crosswise, needs simpler electric conductors and conductors for the means of cutting. In view of the width of the slab, the conveyor duct has relatively short dimensions and requires a smaller intake of energy. The exhauster performance of the gas-exhaust pipe is also less.\nThe transport of the uncut metal slabs, as well as the pieces cut and longitudinally parted, causes no problem, in that two cutting machines are tandem-joined, with a distance between the longitudinal parting direction at least as large as the maximal length of a metal slab.\nA pair of driving rolls is positioned in front of the first cutting machine and a second pair of driving rolls is positioned behind the second cutting machine in longitudinal parting direction. The distance between the first cutting machine and its respective pair of driving rolls is at any time smaller than the uncut length of a minimal metal slab length while the distance between the second cutting machine and its respective pair of driving rolls is smaller than the cut length of a minimal slab after leaving the first cutting machine.\nThe crosswise movability required for the lateral corrections of the cutting machine may be used for an additional function, in that the cutting machine is adjustable crosswise to the longitudinal parting direction to provide for various parting widths of the cut slabs.\nIn addition, the short length of the slag conveyor path advantageously permits the conveyor duct to consist of a vibrating chute. To allow for lateral correction, as well as adjustment for various parting widths, lateral tracing devices are provided for tracing the vertical sides of the metal slabs which determine their width. The cutting machine is arranged to be motionless in the longitudinal parting direction, and the control of the crosswise drive of the cutting arc is geared towards an even distribution of the deviations onto the width of the pieces parted longitudinally.\nWhile the long cutting path on cutting machines which are moved in cutting (transverse) direction require increased input for sound containment, the invention permits the cutting machine, which is stationary in longitudinal parting direction, to be arranged in a soundproof casing corresponding to the size of the cutting machine. The invention also avoids movable attachments and provides a rigid striker bar on at least one side of the roller bed at the level of the metal plate bar thickness.\nA high cutting peformance is attained while using less space when several longitudinal parting paths are arranged parallel next to each other, connected via a device pulling in transversal direction.\nAn example embodiment is shown in the drawing and is described in more detail below."}
{"text": "A Transmission Control Protocol/Internet Protocol (TCP/IP) packet is a relatively commonly used packet. In order to increase an uplink capacity and improve user data transmission efficiency, compression processing needs to be performed on the TCP/IP packet in a transmission process. In the prior art, a header compression manner is used for a TCP/IP packet, that is, compression is performed on a header of the TCP/IP packet. In the header compression manner, if headers of multiple TCP/IP packets vary little and have high repeatability, compression on the headers of the multiple TCP/IP packets in the header compression manner may improve data transmission efficiency and save a network resource.\nIn a data transmission process, in addition to header information of the TCP/IP packet, a packet at an application layer above a TCP/IP layer also has similar content, so that compression may be further performed, so as to increase a transmission capacity. However, compression of application layer data is not implemented in the prior art."}
{"text": "As our daily activities become ever more reliant on access to and interaction with various web services, a growing problem facing a typical end user is that of the ability to streamline and optimize cross-device updates and events originated by web services. A particular challenge is to effectively monitor and navigate to and through the multiple services to which an end user has subscribed to use the desired feature or to access a given update or event. Web with Real Time Communication (WebRTC) services need to notify users and connect sessions, news services alert for news events, sports services alerts for sport scores, stock services alert for stock prices. The assumption of all of these services is that the user has exactly one device that needs to be alerted or, if the user has more than one device, the user will define alerts for each device at the service level. Some notifications need the user to have an active session (e.g. WebRTC). For some devices that get different Internet Protocol (IP) addresses each time they join the network (e.g. laptops), registrations are needed every time a new network address is acquired. Some notifications are sent to persistent addresses (e.g. mailboxes), and others are sent to a device identity whose IP address may change but the network resolves the identity (e.g. smartphones).\nFor many services that require the ability to connect two users together in real-time, the service is based on the user maintaining current registrations with the service.\nToday, most users have multiple networked devices (laptop, smartphone, tablet, etc.,) that may be used at different times. Some of those devices may be shared among several people (e.g. family members) and notifications should come to those devices only when the right individual is using the device. For example, suppose that Grace is a doctor who is on call for consultations from ten separate practices which use separate web-based services to communicate with medical professionals. Each of the services may have a separate unique login and password for Grace. When Grace logs onto the network from any device in the morning, she wants to be connected to all ten services from a single device. If Grace switches from her laptop at home to her cell phone and then to her laptop at work, she may want to inform all ten services to contact her to her work laptop. If Grace has multiple devices registered to all ten services simultaneously, she may want to be able to either have all the devices be informed of incoming messages, or to designate certain devices to communicate with certain services. However, no existing systems provide Grace with a centralized control over the relationship between a plurality of her devices and a plurality of services in the manner described in this example."}
{"text": "To meet rapid progress of industrial development machining techniques also advance constantly to produce finer finished or semi-finished products. The fabrication precision of machinery depends on many factors, such as precision of tools, operation precision of the machinery, fabrication environment, and coupling and holding steadiness of the machinery and tools. One set of machine usually can hold different types of tools alternately to achieve different fabrication results (such as by changing the tools of varying sizes). After each tool changing, a calibration procedure has to be performed (because movements occur during tool changing). This involves a tedious process. Both the chuck and tool have to be removed from the machine and installed again. Fabrication tolerances and errors are bound to happen. To remedy this problem, many techniques have been proposed in prior art. For instance, R.O.C. patent No. M253434 entitled “Detachable chuck” discloses a structure mainly including a cutter stem with an axial conical cavity extended inwards from one end in a tapered manner, a holding member with a conical portion at one end and a coupling portion at the other end and a fastening member coupled on the holding member from outside and screwed on the cutter stem. The conical portion has an inclined peripheral surface formed at an angle the same as that of the conical cavity of the cutter stem, and is held in the conical cavity. The coupling portion has a central hole to hold a cutter under heating. Due to the cutter is held under heating, replacement of the cutter is troublesome. There are still rooms for improvement."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device of the CMOS type, and more particularly relates to a semiconductor device of the CMOS type in which a DRAM element and logic element are formed on the same substrate, and method for production thereof.\n2. Description of the Related Art\nRecently, in seeking to create a more compact and rapid LSI, developments have advanced in dual-mounted DRAM/logic technology in which a DRAM element and logic element are formed on the same substrate. That the DRAM element and logic element are formed on the same substrate has advantages. This is because a large number of wires can be formed between the DRAM element and the logic element within the LSI. In other words, the rate of data transmission can be readily increased, for example graphic processing performance can be improved. For this reason, the field of application is widening into the graphic accelerators for image processing.\nIn general, there are large differences in the production processes for a semiconductor device for a DRAM and a semiconductor device for logic unit. Even when comparing methods of element separation for example, there are differences between those for DRAM elements and those for logic elements.\nNormally, as shown in FIGS. 4 and 5, methods for selective oxidation (referred to as LOCOS 1 and LOCOS 2 herein below) is used in element separation of DRAM elements. Here, LOCOS is an acronym of \"Local Oxidation of Silicon\". First, a brief description of LOCOS 1 will be given with reference to FIG. 4.\nAs shown in FIG. 4(a), a thin oxide film 52 and Si.sub.3 N.sub.4 film 53 are formed, in that order, on the P-type semiconductor substrate 1.\nAs shown in FIG. 4(b), the Si.sub.3 N.sub.4 film 53 and thin oxide film 52 are selectively removed by a widely known technique of photoetching. Then, using the Si.sub.3 N.sub.4 film 52 as a mask, selective oxidation is performed, and a separation oxide film 54 for the purpose of element separation is formed. Following this, the elements are separated and LOCOS-1 completed by the removal of the Si.sub.3 N.sub.4 film 53 and thin oxide film 52, as shown in FIG. 4(c).\nThere are problems in that along with a size-reduction of the DRAM element, when a separation oxide film 54 is formed by LOCOS-1, there are problems in that the interval (pitch) between the element separation region and element formation region cannot be reduced due to the presence of what are known as bird beaks in which the oxide film protrudes in the lateral direction. As a strategy to overcome this problem, and to suppress the spread of oxidation to the element formation region, practitioners have come to use LOCOS-2 in which a polysilicon film is placed below the Si.sub.3 N.sub.4 film and the birds beak absorbed in this area. A brief description of LOCOS-2 will be given using FIG. 5.\nAs shown in FIG. 5(a), a thin oxide film 65, polysilicon film 66, and Si.sub.3 N.sub.4 film 67 are formed, in that order, on a P-type semiconductor substrate 61.\nAs shown in FIG. 5(b), the Si.sub.3 N.sub.4 film 67 is selectively removed by the conventional technique of photoetching, and, using the Si.sub.3 N.sub.4 film 67 as a mask selective oxidation performed, and a separation oxide film 68 is formed. Following this, the separation oxide film is formed as shown in FIG. 5(c), and LOCOS-2 completed by removal of the Si.sub.3 N.sub.4 film 67, polysilicon film 66, and thin oxide film 65.\nMeanwhile, LOCOS-1 and LOCOS-2 have conventionally been used in the element separation of logic elements. Accompanying the change to gate lengths on the half-micron scale, however, has seen the introduction of the use of LOCOS-3. LOCOS-3 is a method in which, following slight preparatory digging of the semiconductor substrate, the area which has been dug is selectively oxidized. The object of LOCOS-3 is to suppress the level-differences which occur as a result of the oxide film. This is because, where there are level-differences occurring in the surface of the substrate in which a separation oxide film has been formed, and where a photoresist acting as a mask has been coated on, the thickness of the photoresist film is not uniform. Where the thickness of the photoresist is non uniform in this way, the thickness of lines of DRAM or other elements formed in the lithography process which follows are not uniform. This is referred to as the standing wave effect. For this reason, LOCOS-3, a method which allows separation oxide films with few level-differences has come into use. A brief description of LOCOS-3 will be given using FIG. 6.\nFirst, as shown in FIG. 6(a), a thin oxide film 79, and Si.sub.3 N.sub.4 film 80 are formed, in that order, on a P-type semiconductor substrate 71.\nAs shown in FIG. 6(b), the Si.sub.3 N.sub.4 film 80, thin oxide film 79, and part of the semiconductor substrate 71 to the required depth, are selectively removed by photoetching. With the Si.sub.3 N.sub.4 film 80 as a mask, selective oxidation is then performed, and a separation oxide film 81 formed. Following this, element separation is carried out as shown in FIG. 6(c), and LOCOS-3 is completed by removal of the Si.sub.3 N.sub.4 film 80 and thin oxide film 79.\nAs described above, in LOCOS-2, LOCOS-2, and LOCOS-3 perform element separation is performed using the separation oxide films from the selective oxidization. Generally however, the LOCOS methods are accompanied by the formation of large level-differences on the surface of the semiconductor substrate as a result of volume expansion which occurs in the selective oxidization. A major feature of LOCOS-3 lies in that there is a preparatory digging out of the silicon substrate by photo-etching to prevent the occurrence of these large level-differences.\nIn each instance of the prior art described above, however, the following drawbacks exist. That is to say, LOCOS-1 and LOCOS-2 are appropriate for the element separation of DRAM elements. Where there is element separation of logic elements formed at a level of about the half-micron scale using LOCOS-1 or LOCOS-2, however, problems arise in that it is difficult to control gate length dimensions. This is because the standing wave effect of in lithography cannot be controlled.\nMoreover, LOCOS-3 is appropriate for the element separation of logic elements. When LOCOS-3 is applied to the element separation of DRAM elements, however, problems arise in that, as a result of the digging out of the semiconductor substrate defects occur in the semiconductor substrate, and diffusion layer leakage increases. Generally, in order to preserve data the diffusion layer leakage is set lower in LSI for DRAM than in LSI for logic. For this reason, the properties of LSI for DRAM deteriorate using the LOCOS-3 method in which the semiconductor substrate is dug out.\nIn addition, LSI for DRAM and LSI for logic have hitherto, been designed and produced using different production techniques. These production techniques are already established. Here, there are different separation oxide films formed by different production techniques. More specifically, the shape and spread of bird beak structures differ. For this reason, when a separation oxide film is formed by one of the methods, (either LOCOS-1 or LOCOS-3 for example), an increase or decrease in the area of the element formation region is induced. The significance of this is that an increase or decrease in the level of LSI integration and/or capacity occurs, and that one of the design assets of LSI for DRAM or logic units already established becomes unusable without further modification.\nWhere LSI with dual-mounted DRAM/logic units are produced, in order to make good use of the respective LSI design assets for DRAM and for logic of the prior art, the design rules for DRAM and for logic elements of the prior art must be used.\nSimilar technology to the present invention is disclosed in Japanese Laid Open Patent Application No.H3-262154. In this LSI different types of separation oxide film are formed on the same substrate. More specifically, it relates to a method for producing a Bi-CMOS type LSI. In this known example, separation oxide films of different thickness are formed on a bipolar transistor and CMOS transistor. Here, the method for forming the separation oxide films is the same."}
{"text": "Heretofore in the serving of casserole dishes or oven heated ware, there have been regular insulated pads known as hot mats which are generally regular in shape such as square or round and containing some type of an insulation material. The mat may be wrapped around the edge of a casserole to hold it when removing it from the oven and carrying it to a table. Additionally, there have in the past been gloves of an insulated type such as barbecue gloves, etc., where the casserole can be gripped for removing it from the oven.\nIn addition, when the hot casserole is placed on the table it is normally desirable to have some type of insulated material between the casserole and the table to prevent burning and discoloration thereof. This in the past has been taken care of by trivets of various types which are elevated above the table and which will allow the hot dish to sit on it without hitting the table.\nHowever, in none of the aforementioned devices is there any provision made for retaining the heat from the casserole during the time it is on the table and the food is served from it.\nIn the past, there have been devices available which do help to retain the heat in the casserole during serving such as the type where a rack is provided to set the casserole in and underneath are candles or other heating elements. Further, there have been provided in the past insulated dishes such as Thermos. Further, there are certain trivets known as electrical trivets which have a mild current passing therethrough.\nHowever, again in none of these devices is there provision for a flexible holder to grip the casserole while it is moved from the oven to the table."}
{"text": "1. Field of the Invention\nThe present invention relates generally to crank and slide presses for metal working. More particularly, the present invention relates to a crank and slide press for redrawing metallic cups into containers, and to a redraw carriage for crank and slide presses that includes improved apparatus for resiliently clamping the redraw cup between the redraw sleeve and a redraw die.\n2. Description of the Prior Art\nCrank and slide presses are used for punching, shearing, drawing, and redrawing operations in manufacturing articles from metallic sheets or rolls of metallic strip material.\nIn general, a crank and slide press includes a crank that is mounted for rotary motion, an electric motor that is connected to the crank and that imparts rotary motion thereto, a slide that is guided for reciprocating motion, and a connecting rod that cooperates with the crank and the slide to change rotary motion into reciprocating motion. A more detailed discussion of crank and slide presses is given by Maytag in U.S. Pat. No. 3,696,657.\nOne specific use for a crank and slide press is in the production of beverage cans. A redraw cup is drawn from a coiled strip by a first crank and slide press; and then the redraw cup is redrawn on a second crank and slide press by forcing the redraw cup through a tool pack that includes a redraw die and a plurality of ironing dies.\nIn the highly competitive container industry, the container must be made with the absolute minimum of material, which means that the finished container must have extremely thin walls. Further, because of the competitiveness of the container industry, the speed of the press must be maximized and downtime of the press must be minimized.\nAs should be apparent, maximizing speed and minimizing downtime are inherently opposite, because higher speeds impose higher stresses on the machinery and cause higher wear rates. Further, it should be apparent that the use of extremely thin walls in the finished containers places stringent accuracy and alignment requirements on such parts as the redraw punch and redraw die.\nIn the redrawing operation, the tool pack, with the redraw die thereof, is mounted circumferentially around the longitudinal axis along which the redraw punch is reciprocated; the redraw punch is reciprocated toward, into, and through the tool pack, including both the redraw die and the ironing dies; and the redraw punch is withdrawn from the tool pack and all of the parts thereof.\nThe redraw cup is positioned against the redraw die and is resiliently held against the redraw die by a redraw sleeve; the redraw sleeve is attached to a redraw carriage; and the redraw carriage is reciprocated toward and away from the redraw die.\nThe redraw sleeve serves two functions. One function of the redraw sleeve is to assure concentricity between the redraw punch and the redraw die. Of course, this concentricity is extremely important because of the thinness of the walls of the finished container and thus the small difference between the diameter of the redraw punch and the redraw die.\nThe other function of the redraw sleeve is to achieve resilient clamping action between a clamping face of the redraw sleeve, the redraw cup, and the redraw die.\nIn the prior art, such as typified by the aforesaid patent to Maytag, the redraw carriage has been actuated toward the redraw die by a cam that is attached to the crank of the press, and by a cam follower mechanism that receives motion from the cam and that transfers this motion to the carriage. It has been customary to use an air cylinder to maintain the cam follower in contact with the cam and to actuate the redraw carriage away from the redraw die.\nIt is important that the redraw sleeve resiliently hold the redraw cup against the redraw die with sufficient force to promote accurate redrawing as opposed to allowing the redraw cup to collapse and wrinkle as it enters the redraw die; but it is equally important to limit the clamping force of the redraw sleeve below a clamping force that will cause any of the redraw cups to rupture rather than being redrawn.\nIn the prior art, as exemplified in the aforesaid patent to Maytag, a spring mechanism has been interposed between the redraw cam and the redraw carriage in order to limit the clamping force of the redraw sleeve. While this construction has been, in general, satisfactory in the past, decreases in the thickness of the walls of beverage containers, and increases in press operating speeds have made the prior art apparatus inadequate for present needs.\nMore particularly, it is well-known that, at high compression velocities, a pressure wave is transmitted through a spring; so that the force actually exerted by the spring varies widely from the static force of the spring. Thus, in redraw operations, the clamping force of a spring will vary considerably from that which the spring would exert under static, or low speed, conditions.\nFurther, it can be appreciated that, with operating speeds of approximately two hundred and fifty redraw operations per minute, the compression wave may be oscillating back and forth in the clamping spring for a higher percentage of the time, or even the entire time, in which the redraw cup is held between the redraw sleeve and the redraw die; so the clamping force of the spring varies greatly during the clamping operation.\nIt can be appreciated that, if the actual clamping force is varying widely due to the compression waves in the clamping spring, then an excessively large portion of the satisfactory range in the clamping force is taken by the spring-caused variations in the clamping force; and the adjustment range for satisfactory operation becomes much more critical.\nFurther, it is important to make adjustments to the redraw press while the press is operating. It can be understood that, if the press is stopped, a small adjustment is made, a short run is made to observe the results of the change, the press is stopped again, etc., considerable production is lost. Because of the loss of production through stopping the press to make adjustments, there is the tendency to avoid making the number of minor adjustments that should be made to optimize the clamping force, and there is a tendency to delay making adjustments as long as possible, even though a more optimum clamping force would result in fewer rejects in the finished containers.\nThus, a principal object of the present invention is to provide apparatus and method for resiliently clamping a redraw cup between a redraw sleeve and a redraw die.\nAnother object of the present invention is to provide apparatus and method for resiliently clamping a redraw cup between a redraw sleeve and a redraw die that provides a uniform pressure during a high percentage of the time in which the redraw cup is clamped.\nStill another object of the present invention is to provide apparatus and method for the clamping operation as described in which the effect of spring pressure-wave forces is eliminated.\nFinally, it is an object of the present invention to provide apparatus and method as described in which the clamping force can be adjusted without stopping the crank and slide press."}
{"text": "The present invention relates generally to computer telephony, and more particularly to an Interactive Voice Response (IVR) system using Voice over Internet Protocol (VoIP) technology in a telecommunications network.\nAn Interactive Voice Response (IVR) system is a software application that accepts a combination of voice telephone input and touch-tone keypad selection and provides appropriate responses. An IVR system is usually a part of a larger application that includes database access.\nAn IVR system uses an embedded software application and has been in commercial use for several years. For example, banks often use an IVR system to allow customers to perform fiscal transactions such as updating the bank account using a telephone or Internet connection. Large businesses routinely use IVR systems in call centers to route incoming calls. Typically in a call center IVR system, to resolve a product issue, a customer dials a customer care telephone number and enters a sequence of touch-tone keypad inputs. After obtaining relevant information regarding the issue, call center IVR system either presents an issue resolution, or simply logs the issue for further investigation. Additionally, movie theaters use IVR systems for selective information lookup such as finding movie schedules, theater locations etc.\nConventional IVRs require specialized architectures to support large applications, and databases containing records for thousands of application users. Consequently IVRs have been implemented in the existing telecommunications network to handle only large applications. Hence, domestic users and small businesses have been unable to enjoy the benefits of an IVR system due to unavailability of a small scale and low cost IVR system.\nThe small scale and IVR system service node was unavailable because the conventional IVRs employed the long-established Public Switched Telephone Network (PSTN) for information transmission. Under the PSTN circuit-switched calls, a communication channel was exclusively dedicated to a conversation for its duration, hence the 64 kbps connection could not be used for any other conversation regardless of whether the caller or the called party talked or remained silent.\nAs a result, the telephone service provider had to bill the calling party for the use of the line for the entire duration of the call. Hence, the IVR systems were too costly for domestic users or small business. Realistically, if a small scale IVR system was introduced under the conventional PSTN system, each domestic user would have blocked a channel of communication for the duration of a call, which was not feasible since the number of potential domestic users far exceeded the number of available communication channels.\nTherefore, a problem is presented in that a personal IVR system must be provided for the benefit of small businesses and domestic users. It is also required that the personal IVR system avoid using the PSTN for signal transmission, since using the PSTN for signal transmission is likely to result in bottlenecking the existing communication channels. What is needed is a node provisioning system capable of offering the benefits of conventional IVR systems to domestic users and small businesses that will use the network resources only when needed. What is needed is an alternative mode of information transmission that can offer the preferred capacity of using the bandwidth of a communication channel only when needed."}
{"text": "Building construction involves precise measurements and calculations in order for a structure to be properly built. One mistake or miscalculation can be magnified many times over and can result in the redoing and recalculation of substandard work. No construction feature requires more precision and accuracy than setting door frames. Typically a door frame is inserted into a rough of a structure being developed. A rough may be defined as a frame opening in a structure including metal or wood studs. Problems associated with the placement of a door frame is the ability of an operator to accurately level, plumb, and square the frame within the rough as it is being attached to the metal or wooden studs. The operator often finds it difficult to hold a frame in place, while making sure the frame is level, plumb and square with one hand, and attaching the frame to the rough with the other hand. Often the frame moves before attachment so that after attachment it is no longer level, plumb or square."}
{"text": "The invention relates to mobile phones, and more particularly, to systems and methods for selectively activating functions provided by a portable device.\nMobile phones are typically equipped with mass storage devices such as flash memory, digital cameras, transmission devices and others, to provide various functions. These embedded devices may be accessed by an electronic device such as a computer, a personal digital assistant (PDA), a programmable consumer electronic device or similar."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of hydrocarbon production from hydrocarbon-bearing formations. More particularly, it concerns fluid compositions that can be useful in improving hydrocarbon production and methods of introducing such fluid compositions into a hydrocarbon-bearing formation.\n2. Description of Related Art\nHydrocarbons (oil, natural gas, etc.) are obtained from a subterranean geologic formation (e.g., a xe2x80x9creservoirxe2x80x9d) by drilling a well that penetrates the hydrocarbon-bearing formation. This provides a partial flowpath for the oil to reach the surface. In order for oil to be xe2x80x9cproduced,xe2x80x9d that is travel from the formation to the well bore (and ultimately to the surface), there must be a sufficiently unimpeded flowpath from the formation to the well bore. Unobstructed flow through the formation rock (e.g., sandstone, carbonates) is possible when rock pores of sufficient size and number are present for the oil to move through the formation.\nA common reason for a decline in oil production is xe2x80x9cdamagexe2x80x9d to the formation that plugs the rock pores and impedes the flow of oil. Often such damage is a result of methods and chemicals used in establishing the well or in remedial operations performed on it. Another reason for lower than expected production is that the formation is naturally xe2x80x9ctightxe2x80x9d (e.g., a low permeability formation), with pores sufficiently small that the oil migrates toward the well bore only very slowly.\nIn general, techniques used to increase the permeability of the formation are referred to as xe2x80x9cstimulation.xe2x80x9d Essentially, one can perform a stimulation technique by: (1) injecting chemicals into the well bore to react with and/or dissolve damage; (2) injecting chemicals through the well bore and into the formation to react with and/or dissolve small portions of the formation to create alternative flowpaths for the hydrocarbon (thus rather than removing the damage, redirecting the migrating oil around the damage); or (3) injecting chemicals through the well bore and into the formation at pressures sufficient to fracture the formation, thereby creating a channel through which hydrocarbon can more readily flow from the formation and into the well bore.\nAcids or acid-based fluids, and fluids containing strong chelants have been found to be useful both for the removal of damage and the creation of alternative flowpaths due to their ability to dissolve certain formation minerals and/or certain contaminants introduced into the well bore and formation during drilling or subsequent operations. The most common agents used in acid treatments of wells are mineral acids, though organic acids are also used. Hydrochloric acid is the preferred acid treatment in carbonate formations. For sandstone formations, the preferred fluid is a hydrochloric/hydrofluoric acid mixture. When the treatment is designed to remove damage resulting from a drilling mud filter cake, it is important to moderate the reaction of the fluid with the filter cake. When a strong acid is delivered in a low-viscosity fluid, often the action of removing the low permeability filter cake is localized, and much of the remainder of the acid is lost to the formation without further removal of the filter cake across the damaged zone. Often drilling mud is designed with a combination of polymers and dissolvable minerals, such as calcium carbonate. An emulsion formulated with an internal acid or internal strong chelant can be used to remove the mineral on a moderated basis by slow release. Dissolution of the mineral compromises cake strength making it more easily removed by production pressures. Acids are also utilized in the process called xe2x80x9cacid fracturingxe2x80x9d for improving permeability in carbonate reservoirs through an etching technique.\nAt present, such acid-based and chelant-based treatments have several serious limitations: lack of uniform radial and axial (along the wellbore) coverage; corrosion of the pumping equipment and well bore tubing; and reaction rates at higher temperatures that are too rapid. An improved well stimulation fluid that has a moderated reactivity and that is less damaging to equipment is a highly desirable goal.\nHydraulic fracturing involves literally breaking or fracturing a portion of the surrounding strata, by injecting a specialized fluid into the well bore directed at the face of the geologic formation at pressures sufficient to initiate and extend a fracture in the formation. Actually, what is created by this process is not always a single fracture, but a fracture zone, that is, a zone in the formation having multiple fractures, through which hydrocarbon can more easily flow to the well bore.\nCertain commonly used fracturing treatments generally comprise at least three principal components: a carrier fluid (usually water or brine), a polymer, and a proppant. Many further comprise a crosslinker. Other compositions used as fracturing fluids include water with additives, viscoelastic surfactant gels, and gelled oils. The purpose of these fracturing fluids is to first create and extend a fracture, and then once it is opened sufficiently, to deliver proppant into the fracture, which keeps the fracture from closing once the pumping operation is completed. The carrier fluid is the means by which proppant and breaker are carried into the formation. A typical fracturing fluid can be prepared by blending a polymer with an aqueous solution (sometimes an oil-based or a multi-phase fluid is desirable); often, the polymer is a solvatable polysaccharide. The purpose of the polymer is to increase the viscosity of the fracturing fluid which aids in the creation.of a fracture; and to thicken the aqueous solution so that solid particles of proppant can be suspended in the solution for delivery into the fracture. In many fracturing treatments, a crosslinking agent is added which further increases the viscosity of the fluid composition by crosslinking the polymer.\nIn addition to being useful in fracturing, fluids that comprise a polymer and crosslinker can also be useful in the workover of a hydrocarbon producing well to improve production. Such workover fluids (e.g., those used in conformance control and zone abandonment operations, among others) with increased viscosity through polymer crosslinking, especially those that become a gel after crosslinking, are useful for isolating particular zones for subsequent treatment by acids, scale removers, or asphaltene diluents, among others. After the treatment, a gel formed by the workover fluid can be intentionally degraded to remove the barrier. In other cases, the gel permanently plugs the zone.\nA chemical or biological agent (e.g., crosslinked polymer, acid, or biocide, among others) that is a useful component of a stimulation or workover fluid can, in certain cases, be neutralized or degraded before reaching the site at which it is to have its effect. Therefore, in certain instances, more of the agent is used in order to be effective and to compensate for agent that is lost in delivering the agent to the site. Thus, there is a need for a more efficient way to deliver useful chemical and biological agents to a desired location in a well.\nThe present invention is directed to fluid compositions that can be used in improving hydrocarbon production and methods of introducing such fluid compositions into a hydrocarbon-bearing formation. The invention relates to delaying the action of at least one chemical or biological agent that is sequestered in the discontinuous phase of an emulsion. Upon exposure to one or more destabilizing conditions, the emulsion is disrupted, releasing the sequestered chemical or biological agent into the bulk fluid of the composition, permitting the agent to have its desired effect.\nOne embodiment is a method of delivering a well treatment fluid composition into a hydrocarbon-bearing formation. In this method, a well treatment fluid that comprises a surfactant stabilized emulsion is pumped into a well. The emulsion comprises a continuous phase and at least one discontinuous phase that comprises at least one chemical or biological agent. The emulsion is destabilized by exposure to at least one of (1) a change in temperature, (2) a change in pH, (3) a change in salinity, (4) a change in alcohol concentration, (5) a change in stabilizing surfactant concentration, (6) a change in organic ion concentration, (7) a change in destabilizing surfactant concentration, (8) a change in surfactant adsorbent material concentration, (9) an ultrasonic pulse, and (10) an electrical field. Destabilization of the emulsion results in the release of the chemical or biological agent from the discontinuous phase.\nAnother embodiment is directed to a method of delivering a well treatment fluid composition, as described above, into a hydrocarbon-bearing formation. In this embodiment, the well treatment fluid is pumped into a well that comprises a first zone and a second zone. The first zone has a first ambient condition, and the second has a second ambient condition that is different from the first ambient condition in at least one respect (e.g., temperature, pH). The fluid composition comes into contact with the first zone, thereby exposing the fluid composition to the first ambient condition, and later the fluid composition comes into contact with the second zone, exposing it to the second ambient condition. The emulsion of the well treatment fluid remains stable for at least about one minute when the fluid is in contact with the ambient condition in the first zone. However contact with the ambient condition in the second zone destabilizes the emulsion, thereby releasing the chemical or biological agent. The ambient conditions of the first zone and the second zone differ in at least one of (1) temperature, (2) pH, (3) salinity, (4) alcohol concentration, (5) stabilizing surfactant concentration, (6) organic ion concentration, (7) destabilizing surfactant concentration, (8) surfactant adsorbent material concentration, (9) ultrasonic energy, and (10) an electrical field.\nYet another embodiment is directed to a method of delivering a well treatment fluid composition into a hydrocarbon-bearing formation, in which, prior to pumping the fluid into the well, at least one chemical that undergoes a chemical reaction or that dissolves in the well treatment fluid composition is added and mixed into the composition. As the chemical reaction of the chemical proceeds or as it dissolves, it causes a change in at least one of the following characteristics of the well treatment fluid: (1) pH, (2) alcohol concentration, (3) stabilizing surfactant concentration, (4) organic ion concentration, (5) salinity, and (6) destabilizing surfactant concentration. The change in the fluid that results is sufficient to destabilize the emulsion, thus causing the chemical or biological agent to be released. Naturally the pumping of this well treatment fluid must occur before the reaction of the chemical or dissolution of the chemical goes to completion; if this were not the case the release of the chemical or biological agent would not be sufficiently delayed.\nStill another embodiment is directed to a well treatment fluid composition, comprising a surfactant stabilized emulsion, wherein the emulsion comprises a continuous phase and at least one discontinuous phase, the discontinuous phase comprising at least one chemical or biological agent. The emulsion in the fluid composition will destabilize sufficiently to release the chemical or biological agent upon exposure to at least one of the conditions listed above. Well treatment fluid compositions of the present invention can be used, for example, in stimulation and workover of oil, gas and water wells, and thus these compositions can be fracturing fluids, acidizing treatment fluids, zone abandonment fluids, conformance control fluids or drilling mud filter cake cleanup fluids depending on their composition and the manner of their application. Well treatment-fluids of the present invention can also be biocide treatment fluids, corrosion inhibition fluids, or breaker fluids.\nIt is preferred that the emulsion of the fluid compositions of the present invention be a microemulsion, because microemulsions tend to have greater stability than macroemulsions. Preferably, the chemical or biological agent initially in the discontinuous phase is selected from the group consisting of crosslinkers, polymers, biocides, corrosion inhibitors, corrosion dissolvers, pH modifiers, breakers, metal chelators, metal complexors, antioxidants, wetting agents, polymer stabilizers, clay stabilizers, scale inhibitors, scale dissolvers, wax inhibitors, wax dissolvers, asphaltene precipitation inhibitors, waterflow inhibitors, sand consolidation chemicals permeability modifiers, foaming agents, microorganisms, nutrients for microorganisms, and salts. Fluid compositions and methods of the present invention can be used for improved fracturing or workover operations by permitting a delay of crosslinking. For example, a preferred embodiment of the present invention is a fluid composition that comprises a cross-linkable polymer in the continuous phase of the emulsion and a crosslinker in the discontinuous phase of the emulsion. Eventually, the crosslinker is released into the continuous phase due to exposure of the emulsion to an emulsion destabilizing effect. Thus, the crosslinking of the polymer is delayed until the emulsion is destabilized. Alternatively, instead of sequestering the crosslinker to delay crosslinking, the polymer or a pH modifier (the pH modifier causes the pH of the fluid composition to change to a pH favorable for the crosslinking reaction) can be sequestered in the discontinuous phase of an emulsion until the emulsion is disrupted. Thus, in certain embodiments, it is preferred that the polymer and the crosslinker can react with each other readily when the emulsion is destabilized. Such fluid compositions comprising a polymer that becomes crosslinked are particularly useful as fracturing fluids, but can also be used in workover operations.\nThe delay in crosslinking afforded by certain embodiments of the present invention is beneficial in that the amount of energy required to pump the fluids can be reduced, the penetration of certain fluids can be improved, and shear and frictional damage to polymers can be reduced. By delaying crosslinking, crosslinkers can be more thoroughly mixed with the polymer fluid prior to crosslink initiation, providing more effective crosslinks, more uniform distribution of crosslinks, and better gel properties. Certain embodiments provide better control of crosslinking than certain mechanisms known in the art (e.g., use of delaying ligand for metal crosslinkers).\nFurthermore, certain embodiments of the present invention can be used to overcome operational problems at the surface. For example, a wetting surfactant which typically causes problems in pumping due to foaming can be sequestered in the discontinuous phase of the emulsion of a fluid composition of the present invention until its release is triggered by a destabilizing condition downhole.\nDelayed release of certain chemicals in various stimulation fluids and workover fluids is also desirable for other reasons. The use of certain methods and compositions of the present invention in hydrocarbon well operations can reduce the damage sustained by tubulars and pumping equipment. For example, mineral acids (e.g., pH modifiers) used in acidizing treatments are corrosive and can cause damage to the tubulars and pumping equipment. If the acid is sequestered and its release delayed until it is in the formation or shortly before it reaches the desired location downhole, damage to pumping equipment and tubulars can be reduced or avoided. Thus, certain embodiments of the present invention are directed to the delayed release of chemicals to minimize damage to pumping equipment and/or well bore structural elements.\nFurthermore, using certain methods and compositions of the present invention, improved radial penetration of the hydrocarbon containing formation by certain workover fluids can be achieved. In certain instances, chemical or biological agents, that would normally react in the immediate vicinity of the well bore, are sequestered for a time in an emulsion, allowing them to travel further radially through the formation before reacting.\nAnother advantage of the present invention is that certain embodiments can be used for slow release of corrosion inhibitors and dissolvers and biocides downhole. Corrosion inhibitors can be slowly released in the well bore to enhance their functionality over time. Biocides can be slowly released to achieve protection over a longer time span, thus requiring lower overall quantities of the biocide."}
{"text": "The present invention relates to computer systems, and more particularly to a method and system for improving performance of a RAID controller through adaptive write back/write through caching.\nRedundant array of inexpensive disks (xe2x80x9cRAIDxe2x80x9d) technology is increasingly used in computer systems to provide increased capacity and performance. FIG. 1 depicts a conventional RAID subsystem 20 that is coupled to a computer system 10. The computer system 10 is typically a server which allows a large number of users to access the RAID subsystem 20. The RAID subsystem 20 includes a RAID controller 22 having a cache 24. The RAID subsystem 20 also includes disks 26, 28, 30 and 32. For clarity, only these portions of the RAID subsystem 20 are depicted.\nThe RAID controller 22 provides an interface between the RAID subsystem 20 and the computer system 10. Thus, the RAID controller 22 includes the hardware that interfaces between the computer system 10 and the disks 26, 28, 30 and 32 and allows the disks 26, 28, 30 and 32 to appear to the computer system 10 as a single disk. The RAID controller 22 also includes the cache 24. The cache is used to improve input to and output from the RAID subsystem 20 by reducing the latency and increasing throughput. This is possible because cache accesses are typically faster than accesses of a disk 26, 28, 30 or 32. When data is to be output from the RAID subsystem 20, data can be written from the disks 26, 28, 30 and 32 to the cache 24. The data can then be output from the cache 24, typically at a higher rate than data can be written from the disks 26, 28, 30 and 32. Performance is improved when data written to the cache by the disk is repeatedly accessed. When this happens disk I/O does not occur and performance is improved. When data is to be written to the RAID subsystem 20, data can be input to the cache 24. The data can then be asynchronously written to the disks 26, 28, 30 and 32, typically at a lower rate than data is written to the cache 24.\nTypically, the cache 24 can be in one of two configurations. The cache 24 can be configured so that data stored in the cache is written to the RAID subsystem 20 at the same time it is written to the cache. This configuration is known as write through. In write through mode, data written to the RAID subsystem 20 is written to both the cache 24 and to the disks 26, 28, 30 and 32. This may improve performance because if data in the cache 24 is reused it can be retrieved from the cache 24 without the delay of a disk access. In the case of write through even though data is written to the cache 24 the operation does not complete until the data is actually written to a disk 26, 28, 30, or 32. Alternatively, the cache 24 can be placed in write back mode. In write back, data is written directly to the cache 24 without requiring data be immediately written to the disks 26, 28, 30 or 32. This allows a faster completion of the write operation because a write operation to the cache 24 is much faster than a write operation to the disks 26, 28, 30 or 32. Thus, the operation need not wait for the disk write to complete. At a later time, the modified data in the cache 24 is written from the cache 24 to the disks 26, 28, 30 and 32 of the RAID subsystem 20.\nIn general, the cache 24 is placed in write back mode. Thus, the cache 24 is used to store data that is written to the disks 26, 28, 30 and 32 without first writing the data to disk. In write back mode the cache is also used to store data that is read from the disk. In many instances, write back mode improves the performance of the RAID subsystem 20, as described above. However, one of ordinary skill in the art will readily recognize that the write back mode may not provide optimal performance in some instances. In particular, when the cache 24 starts to become full, an I/O operation may stall until a portion of the cache 24 is flushed. As a result, performance of the RAID subsystem 20 suffers because any outstanding disk write command must wait for a cache buffer management process in the RAID controller 22 to flush a modified buffer (not explicitly shown in FIG. 1) and make the cache 24 available for the new write operation. Write through operations also require a buffer of the cache 24 to store the write data, but the buffer management does not have to flush modified data to free a buffer of the cache 24 because all buffer data is already written to the disk 26, 28, 30 and 32. In this case buffer management can simply overwrite any existing buffer in the cache 24 to make space for an outstanding write command. In this case, the write through operation may complete without delay.\nOne solution to performance problems due to the cache 24 becoming full is to make the cache 24 larger. A larger cache 24 will become full less frequently. As a result, it was theorized that performance of the RAID subsystem 20 should improve. However, the increase in the size of the memory of the cache 24 often has the opposite result. An increase in the size of the memory of the cache 24 can actually degrade the performance of the RAID subsystem 20. This occurs because if a larger cache becomes full it will require more Raid Controller 22 buffer management-processing overhead to manage and free buffers. This buffer management overhead degrades performance. This condition occurs whenever the workload consists of a sustained I/O rate that is greater than the aggregate capabilities of the RAID subsystem 20.\nAccordingly, what is needed is a system and method for improving the performance of a RAID subsystem 20 that includes a cache 24. The present invention addresses such a need.\nThe present invention provides a method and system for accessing data in a redundant array of inexpensive disks (RAID) subsystem. The RAID subsystem includes a RAID controller having a cache and a plurality of disks. The method and system comprise utilizing the cache in a write back mode if the RAID subsystem is lightly loaded. In write back mode, the data is written to the cache prior to storing the data on at least one disk of the plurality of disks or prior to outputting the data from the RAID subsystem. The method and system also comprise utilizing the cache in a write through mode if the RAID subsystem is heavily loaded. In the write through mode, the data is written directly to at least one disk of the plurality of disks.\nAccording to the system and method disclosed herein, the present invention provides a faster, more efficient mechanism for accessing data in a RAID subsystem."}
{"text": "The art prior to the invention shown in my copending application Ser. No. 06/502,763 filed June 9, 1983, did not provide an acceptable locking/quick release socket wrench extension.\nApplication Ser. No. 06/502,763 discloses a quick release locking arrangement for a socket wrench extension having a square drive end in which a transverse bore retains lock means for movement between a lock position in which the locking means projects from one end of the bore far enough to engage a cooperating cavity or undercut in the drive recess of a wrench socket (or other tool) and a release position in which the lock means lies within the cross section of the drive end and permits removal of the drive end from the socket. A lock tab is mounted for movement between a first position obstructing movement of the lock means to its release position and to a second position in which the lock means can move to the release position. The lock tab is a part of and controlled by a sleeve which is biased to hold the lock tab in the first position.\nMy prior invention required manual actuation of the sleeve to mount a socket on the extension. The design did not have \"push-on\" mounting capability."}
{"text": "Fuel cells are well-known and are commonly used to produce electrical energy from reducing and oxidizing reactant fluids to power electrical apparatus, such as apparatus on-board space vehicles, transportation vehicles, or as on-site generators for buildings. Each individual fuel cell generally includes an anode catalyst and a cathode catalyst separated by an electrolyte, such as a proton exchange membrane (“PEM”) as known in the art. Frequently, a diffusion layer is secured between the catalyst and a substrate layer. The substrate layer is usually secured between the diffusion layer and a flow field. If there is no diffusion layer, the substrate is secured between the catalyst and the flow field. Flow fields define flow channels for directing reactant streams through the fuel cell in fluid communication through the diffusion and substrate layers with the catalysts. As is known, the flow fields may be porous water transport plates or solid separator plates.\nThe diffusion layer is typically a highly porous, electrical conductor made from carbon black and a hydrophobic polymer, such as polytetrafluoroethylene. Depending on the fuel cell design, the diffusion layer is usually hydrophobic, however, the diffusion layer may be partially hydrophobic and partially hydrophilic to facilitate simultaneous liquid and gaseous transport through the layer. The diffusion layer is usually about 25-100 microns thick. The diffusion layer facilitates transfer of the reactant streams through the fuel cell by minimizing the thickness of water films on the surface of the catalysts. The diffusion layer also facilitates the removal of product water from the fuel cell.\nThe substrate layer adjacent the diffusion layer is highly porous and made from expensive carbon fibers and a well known manufacturing process that requires high temperature graphitizing. Depending on the cell design, the substrate layer may be either hydrophobic or hydrophilic. The substrate is usually about 150-300 microns thick. The substrate facilitates the transport of reactant streams, water vapor, liquid water and electrons. The substrate conducts electrons both through the plane of the substrate and in the plane of the substrate from a centerline of an adjacent flow channel to ribs of the flow field, wherein the flow channel is defined between ribs of the flow field. The substrate also facilitates the diffusion and flow of reactant streams and product water both through the plane of the substrate and in the plane of the substrate from the flow channel to the centerline of the flow channel beneath the rib of the flow field, or vice versa. Also, the substrate must have a flexural strength adequate to distribute an axial pressure load relatively uniformly over total surface areas of adjacent layers. For example, the pressure load must be distributed evenly across the flow channels of the flow field to prevent the substrate from deforming into the flow channels. An exemplary flexural strength is about 200 kilogram force per square centimeter (“kgf/cm2”). U.S. Pat. No. 4,851,304 to Miwa et al. describes properties of typical fuel cell substrates. Planar types of fuel cells are secured in compression in a fuel cell stack by a combination of pressure plates and tie-rods, as is well known. This axial compressive force minimizes the resistance of the cells and is required to obtain suitable fluid seals. The substrate must have a compressive strength that is typically at least two times the axial force on the cell stack. Typical compressive strength of a prior art substrate is greater than 10 kgf/cm2.\nIn fuel cells of the prior art, it is known that carbon or graphite layers or papers are traditionally secured between the catalysts and flow fields as diffusion and/or substrate layers. However, the use of known carbon or graphite layers or papers presents significant problems, including high manufacturing costs, impeding the diffusion of hydrogen and oxygen through pores defined by the layers, and impeding the outflow of fuel cell product water from a cathode catalyst.\nA partial solution to the use of carbon or graphite paper is disclosed in U.S. Pat. No. 5,707,755, entitled “PEM/SPE Fuel Cell” that issued on Jan. 13, 1998 to Grot. The patent discloses, instead of the carbon or graphite layers or papers, the use of a plurality of electrically conductive filaments secured with a specific orientation with respect to flow channels, or grooves of a flow field. The specific orientation is longitudinal so that the filaments extend across, and do not fall into, the flow channels, or grooves, that direct reactant flow through the fuel cell.\nNonetheless, the prior art has limitations. Manufacturing a fuel cell with conductive filaments having such a specific orientation imposes substantial cost and manufacturing burdens. Moreover, the prior art fails to rectify the significant substrate manufacturing costs associated with the mass production of fuel cells. The substrates are costly because they include expensive carbon fibers and are manufactured through a costly high temperature graphitizing process well known in the art. Accordingly, there is a need for a fuel cell that minimizes substrate costs by replacing the prior art substrates with a cost effective material."}
{"text": "Mobile devices e.g., telephones and PDA's, are commonly used to send and receive communications, such as electronic mail, text messages, voice messages, and telephone calls. Mobile devices provide the advantage of convenience and allow users to access communications from any location where there is a signal. However, the small screen sizes and storage capacity of mobile devices make it difficult for users to view, sort, organize, and manipulate communications from a mobile device. Current applications do not provide convenient ways for users to view communications compactly and perform actions on groups of communications efficiently.\nIt is with respect to these and other considerations that embodiments of the present invention have been made. Also, although relatively specific problems have been discussed, it should be understood that embodiments of the present invention should not be limited to solving the specific problems identified in the background."}
{"text": "As technology advances, there is a need to combine multiple chips into a device. Combining chips into a device, for example, requires an interposer with larger die size so as to allow more chips to be disposed on the interposer. The size of the die for use as an interposer, however, is limited by lithography tool hardware.\nThere is, therefore, a need to provide effective stitching of dies to reduce footprint, improve performance and reduce cost."}
{"text": "A suction filter of this type is described in German Patent No. DE-PS 35 33 038. A cage is placed on the housing, through which the agitator shaft is led and on which the drive motor is placed. For the vertical adjustment of the agitator shaft two lateral piston rods are provided, which are connected to each other by means of a main cross beam. In turn, two hydraulic cylinders projecting vertically upward are arranged on the main cross beam itself. The piston rods of these hydraulic cylinders are connected to each other by means of a cross beam.\nA disadvantage of this embodiment is, however, that such an agitator suction filter has a relatively large overall height."}
{"text": "The present invention relates to a receiver for a coded electronic security system, and more particularly, to a receiver capable of monitoring a plurality of different transmitters and indicating which of the transmitters has been activated by an alarm condition.\nSecurity systems for protecting homes and businesses are becoming increasingly popular due to an increase in vandalism and theft. The most commonly used protective systems require wiring of doors and windows in such a manner that the unauthorized opening of a protected door or window activates an alarm. ConventIonal systems which require extensive wiring and specially designed switching devices are susceptible to tampering and failure, and they are expensive to install.\nHeretofore, attempts at using wireless intrusion alarm systems have met with limited success because such systems have been simply designed merely to indicate the presence of an actuating signal at the remote receiver. Any discrimination relating to rejection of other signals also present was provided merely to avoid false alarms. An indication of an alarm generally did not provide an indicatIon of which particular transmitter in a system having a plurality of transmitters had activated the alarm, unless there were a like number of transmitters and receivers each operating on a different frequency within the system."}
{"text": "Automotive vehicle A/C systems require maintenance. Maintenance includes refrigerant recovery, evacuation and recharging of the A/C system. Portable refrigerant recovery units are used in connection with this maintenance. The refrigerant recovery units connect to the A/C system of the automotive vehicle to recover refrigerant out of the system, separate out contaminants and oil, and recharge the system with additional refrigerant. During service it is desirable to recover, recycle or reclaim the refrigerant."}
{"text": "Parkinson's disease (PD) is a progressive neurodegenerative disorder that is characterized phenotypically by bradykinesia, rigidity, tremor, and neuropsychiatric disturbances (Savitt et al., 2006). Although the cause of PD in the majority of cases is unknown, there are rare familial cases for which the genes have been identified. There are at least sixteen PD associated loci (Gasser, 2007). Mutations in α-synuclein and leucine rich repeat kinase 2 (LRRK2) cause autosomal dominant PD. Four genes have been linked to autosomal recessive PD (AR-PD) and include mutations in parkin, DJ-1, PINK1, and ATP13A2. Investigating the biology of these genes and their mutant protein has provided tremendous insight into the pathogenesis of both familial and sporadic PD (Gasser, 2007; Savitt et al., 2006).\nParkin is an ubiquitin E3 ligase (Shimura et al., 2000; Zhang et al., 2000). In general, PD-associated mutations in parkin lead to loss of its E3 ligase function (Tanaka et al., 2004). Moreover, oxidative, nitrosative, and dopaminergic stress, which play important pathogenic roles in PD, inactivate parkin, suggesting that parkin inactivation may play a role in sporadic PD (Chung et al., 2004; LaVoie et al., 2005; Winklhofer et al., 2003). Thus, substrates of parkin that are subject to proteasomal degradation should accumulate in animal and cellular models of parkin inactivation and AR-PD due to parkin mutations, and also in sporadic PD. There are a diverse array of parkin substrates that has hindered the generation of a consensus in the field on parkin's physiologic function and pathologic role in PD. Moreover, parkin's ability to mono- and poly-ubiquitinate, as well as, ubiquitinate proteins with both lysine-48 and lysine-63 chains has made it difficult to reconcile a common biochemical pathway for parkin's role in PD (Dawson and Dawson, 2010).\nA new Parkin Interacting Substrate, PARIS, has been identified, which provides a molecular mechanism for neurodegeneration due to parkin inactivation in PD. Parkin regulates the levels of PARIS via the ubiquitin proteasome system (UPS). PARIS is a major transcriptional repressor of peroxisome proliferator-activated receptor gamma (PPARγ) coactivator-1α (PGC-1α) expression and that conditional knockout (KO) of parkin in adult mice leads to progressive loss of dopamine (DA) neurons through PARIS overexpression and transcriptional repression of PGC-1α."}
{"text": "The escalating demands for high density and performance associated with non-volatile memory devices require small design features, high reliability and increased manufacturing throughput. The reduction of design features, however, challenges the limitations of conventional methodology. For example, the reduction of design features makes it difficult for the memory device to meet its expected data retention requirement.\nIn addition, various processes that are used to complete the formation of a working memory device, such as an electrically erasable programmable read only memory (EEPROM) device, often create problems associated with the operation of the memory device. For example, it is often difficult to deposit an interlayer dielectric that fills all the space between adjacent memory cells. This often leads to gaps or voids in the interlayer dielectric. Such gaps or voids in the interlayer dielectric may lead to charge leakage problems associated with memory cells in the memory device. These voids may also make it difficult to program and/or erase the memory device in an efficient manner and, ultimately, may lead to device failure."}
{"text": "The present invention relates to headlights for self-propelled vehicles, especially for motor vehicles.\nA headlight for a self-powered vehicle is described in German Published Patent Application 24 31 095 A1. This headlight has a reflector, which is pivotally mounted in a holder and which is adjustable by means of an adjusting mechanism. The adjusting mechanism has an eccentric adjusting element engaging eccentrically to a pivot axis of the reflector, which is connected with the reflector by a universal joint. This pivot joint is formed by a pin which has two ball ends so that the pin is pivotally engaged with one ball end in a spherical receptacle connected with the reflector and is engaged with the other ball end in another spherical receptacle provided with a threaded passage. A threaded portion of the adjusting element is screwed in the threaded passage. This double pivotal connection is required to minimize bending moments and/or transverse forces of the adjusting element, which can occur because of the pivoting motion of the reflector caused by the adjusting element. The known double-jointed pivot joint with the pin and both spherical receptacles requires comparatively many parts and takes up a comparatively large space, particularly in the adjusting direction of the adjusting element, because it requires the pin provided with the ball ends. This double-jointed pivot joint can be used only when there is enough room."}
{"text": "Cerebral embolism is a known complication of cardiac surgery, cardiopulmonary bypass and catheter-based interventional cardiology and electrophysiology procedures. Embolic particles, which may include thrombus, atheroma and lipids, may become dislodged by surgical or catheter manipulations and enter the bloodstream, embolizing in the brain or other vital organs downstream. Cerebral embolism can lead to neuropsychological deficits, stroke and even death. Prevention of cerebral embolism benefits patients and improves the outcome of these procedures.\nVarious embolic protection devices are known in the art. An embolic protection device for side branch vessels of the aortic arch has for instance been disclosed in US 2004/0215167. This embolic protection device has an expandable tubular structure supporting a filter mesh material. The embolic protection device is compressed to a small diameter for insertion into a patient's aorta, and then expanded within the aorta with the filter mesh material positioned to allow blood to enter side branch vessels connected to the aorta and to prevent embolic material from entering the side branch vessels. The device is deployed and left in place for long-term protection. Alternatively, the device may be compressed and withdrawn from the aorta.\nHowever, the embolic protection devices disclosed in 2004/0215167 have a number of drawbacks. The device may be difficult to extract from the aortic arch as a stent like design is devised for permanent implantation and removing a stent may harm the implantation site. The device also forms along the aorta and may at least partly be pressed against or into the ostia regions of the side vessels. Most often these ostia regions are subject to sedimented plaque on the outside of the tissue in these ostia regions. When a stent like device is pressed against the plaque, the latter loosens from the tissue on which it is situated and is washed along the side branch vessels as debris. However, this debris is an undesired embolic material, which the device should avoid to enter the branch vessels.\nIn U.S. Pat. No. 6,258,120 implantable cerebral protection device is disclosed for diverting emboli away from the carotid arteries in the aorta. The disclosed devices are aortic diverters that generally comprise a hollow tube with a substantially cylindrical or conical wall, which is impermeable to emboli and which has open ends that allow blood to enter one end, flow through the tube and exit the other end. The proximal end of the hollow tube is circumferentially sized to completely fill the lumen of the aorta. Additionally, snowshoe aortic diverters, which are planar rather than cylindrical, are also disclosed. The methods disclosed in U.S. Pat. No. 6,258,120 include the steps of providing an aortic diverter carried by an intravascular catheter, introducing the intravascular catheter into the vascular system, advancing the intravascular catheter into the aortic arch to the region of the carotid arteries, and deploying the aortic diverter.\nHowever, like the devices disclosed in US 2004/0215167, the devices and methods of U.S. Pat. No. 6,258,120 may damage the aortic vessel wall. Furthermore, a leakage of embolic material into the side branch vessels of the aortic arch may be present, e.g. past the periphery of the tubular structure, or of the snowshoe like embodiments disclosed. Moreover, the devices disclosed may, at least partly, contact the ostia of the side branch vessels, and thus set free embolic debris from the ostia which is carried to the carotid arteries and may lead to cerebral damage. Further, a backflow carrying embolic material may occur from the distal end of the devices of U.S. Pat. No. 6,258,120 into the side branch vessels. The snowshoe like devices have and attached handle or cannula and need to be installed by means of open chest surgery comprising incising the aorta, which has numerous drawbacks compared to intravascular delivery, including aortic trauma. The devices need to be secured to the lumen of the aorta through various mechanisms including sutures, surgical clips, hooks, adhesive material, substantially rigid sleeves, or frictional engagement. Such securing is difficult to accomplish in a reliable manner via transvascular access.\nIn US 2008/0065145 an embolic protection device and method of use are disclosed. A blood debris deflector umbrella is disclosed, which has a blood flow permeable covering. The umbrella is extending over the ostia of the brachiocephalic artery and the left carotid artery. The deflector is inserted percutaneously and placed by means of a catheter, either via the right arteria subclavia ending in the aortic arch via the brachiocephalic artery, or the femoral artery and into the brachiocephalic artery and the right arteria subclavia. However, the device has in any case a guide wire arranged extending between the aortic arch and the brachiocephalic artery, which is to be protected by the device. This means that, like in the aforementioned disclosures, the device will inevitably be intravascularly manipulated and thus it will likely contact the ostia of the brachiocephalic artery, i.e. the side branch vessel of the aorta leading to the right carotid artery. Thus a risk for iatrogenic caused embolization is present, i.e. the physician likely will set free embolic debris from the ostia of the brachiocephalic artery when using the device. The embolic debris is carried to the right carotid artery and may lead to cerebral damage.\nMoreover, the dome shaped device of US 2008/0065145 appears to be difficult, to work in practice due to the anatomical structure and position of the afore described access way via the brachiocephalic artery. The device will have to be very large in order to cover the ostia of the brachiocephalic artery and the left carotid artery. Thus the device will be very voluminous in the aortic arch.\nHence, known cerebral embolic protection devices have shortcomings, including: difficult to position in a vessel, and even more difficult to position in two vessels; they may cause damage to the vessel wall and potentially cause an emboli themselves; they are hindering surgeons when trying to achieve a good result with some intended intervention/operation; visualization of the protective device may impair visualization of other components used during concurrent medical procedures; they may cause impaired flow if they are designed to collect the embolic material.\nThus, there is a need for a new, or improved, or alternative device, or method for preventing embolic material from entering branch vessels, such as the aortic arch side branch vessels, and/or from creating debris from the ostia of the aortic arch side branch vessels that may be carried towards the brain of a patient during a medical procedure.\nHence, an improved embolic protection device or method would be advantageous."}
{"text": "Not applicable.\nNot applicable.\nThe present invention relates generally to the field of bicycles, and more specifically, to control systems such as brake, gear and transmission, and other systems employing wire and cables for use therewith.\nBraking systems used for cycles and other vehicles, such as bicycles typically comprise two principle types. One type, known as a coaster brake, includes braking components conventionally housed in the hub of the rear wheel of the bicycle and is operated by pressing on the bicycle pedals in a direction rotationally opposite to the forward pedaling direction to achieve the desired braking operation. The other type of bicycle brake system is known as a caliper brake. Caliper brakes are comprised of dual opposing caliper levers connected for cooperative pivoting motion to the bicycle frame to encompass corresponding opposing circumferential portions of the wheel rim to which a tire is mounted, whether front-wheel mounted or rear-wheel mounted, so that the brake pad-lined ends of the opposing caliper levers are spaced apart from corresponding sides of the wheel rim to which they are closely positioned in the un-braked position. The lever system of the caliper brake operates when the caliper lever ends are squeezed together into a braking position by a brake cable that is pulled when the bicycle rider squeezes an operatively-connected brake hand lever mounted on the handlebars of the bicycle. The squeezing motion causes the brakes to be applied by providing frictional contact between the brake pads levers and the front and/or rear wheel rims according to the installed braking configuration, and hence the desired braking performance depending on the pulling force exerted by the rider. Disc brake systems are also known to be used with bicycles. No matter the type of braking system, cabling system requirements are substantially similar in that a tension is required to actuate the bakes, provided through the cable.\nThe above-described brakes are thus considered to provide an efficient motion-retarding mechanism for slowing and stopping a bicycle. However, the traditional placement of the brake hand levers about the handlebars of the bicycle may not be optimal for ready, convenient operation when the rider is astride a mountain bike or a street bike (on-road or off-road), whether mountain biking or traveling over long distances because of the differing positions adopted by the rider depending on terrain and slope. In particular, the rider may adopt a standing, hunched-over, or semi-standing position when riding uphill, downhill or over rugged terrain, and thus must retain his balance by altering his grip on different portions of the handlebars and in differing positions when gripping those portions in contrast to the conventional sitting or cruising position to which hand brake levers are commonly oriented. Thus, an important limitation in the related art is the failure to accommodate the rider by providing direct access to appropriately positioned brake hand levers in the most significant of those varying riding positions, with commensurate compromises in rider control in those different positions not directly adjacent to the conventionally-placed hand brake levers.\nTherefore, it is desirable to provide an improved cable-handling system for use with cycles that overcomes these prior art shortcomings.\nThe present invention provides an important improvement in the bicycle-brake system art, by providing a supplemental apparatus for operationally connecting an auxiliary bicycle brake actuator apparatus to an existing bicycle brake system, thereby overcoming the prior art shortcomings described above. More specifically, the invention includes a cable joining system for adding additional manual interface(s), i.e. brake hand lever with connecting elements, for use with an existing braking device and system and thus providing enhanced brake system operability by the bicycle rider in additional riding positions, such as those intermediate positions commonly adopted by street bike riders and mountain bike riders. This improvement is readily achieved by retrofitting the original braking system with the supplemental braking apparatus, without cutting or otherwise destructively modifying the existing cable system provided on the bicycle. Accordingly, the supplemental brake lever mechanism may be positioned to be readily operated by the rider in an alternate riding position not previously amenable to direct operation of a selected front wheel or rear wheel brake system.\nAs installed, the inventive system is secured to the frame of the cycle (such as a bicycle, tricycle, quadricycle, scooter, which may be pedal-powered, motor assisted or a hybrid thereof) having an existing brake system with forward and rear cable stops provided at spaced intervals along the top tube or other structural element of the cycle. Such structural element may further include a handlebar, fork, seat post or even the fork structure, i.e., any member other than a rotational member such as a wheel. Accordingly, such structural element is utilized to secure one or more brake wires (and optionally, gear shifting wires) therealong, and it will be understood that all references to xe2x80x9cbrakexe2x80x9d wires will encompass control wires for gearing and transmission systems or other systems for use with the cycle. It will be further understood that further references to xe2x80x9cbicyclexe2x80x9d hereafter includes all of the above-listed species of vehicles including cycles. According to the invention, an auxiliary brake handle (actuator) is secured to the handlebars of the bicycle and a newly installed auxiliary brake wire, connected to the auxiliary brake handle, is secured to the main brake wire. According to the invention, the newly installed auxiliary brake wire and main brake wire are assembled and secured in a xe2x80x9cYxe2x80x9d configuration between the forward and rear cable stops, with a clamp serving as a joining node for securing the auxiliary brake wire to the main brake wire. Accordingly, the present invention causes activation of the auxiliary brake wire, now directly secured to and linked with the main brake wire, upon an actuation of the auxiliary brake system, to subsequently cause a pull upon the main brake wire to initiate braking operation of the operatively connected brake calipers and provide the desired braking action of the corresponding wheel, while also preserving operability of the brake system by the main brake handle.\nThus, according to the invention, the cable joining system of the present invention allows a user to add yet additional manual interfaces for additional cable-based devices (brakes, gearing, etc.) without cutting or modifying the existing cable system or operability of the calipers or braking performance. Indeed, the invention provides enhanced brake system operability by virtue of the rider\"\"s more direct control and proximity to the auxiliary brake handle in riding positions different from the singular position (cruising) to which most brake systems are installed. The inventive system comprises four elements: 1) a housing joiner, 2) an interconnecting auxiliary brake cable or wire, 3) a clamp for joining the interconnecting brake cable or wire any cycle-mounted operational cable or wire, and 4) lever or actuator for actuating the auxiliary brake cable or wire. The housing joiner and clamp (such as a bullet clamp) cooperatively function with the existing cable stops provided on the bicycle frame, so that the two manual interfaces may independently operate the calipers of the selected brake system.\nThe invention thus provides a significant advancement in the art, as it is a simple method of adding an after-market braking system (or gearing system or other cable/wire actuated system) interface to an existing bicycle braking system utilizing a brake cable system including a brake wire sheathed in a cable housing. The method of system installation of the present invention is achieved by squeezing together the calipers of the selected brake system to release the tension of the main brake wire and housing (sheath), thereby enabling unhooking of the cable housing from the cable stops while retaining the brake wire in the installed adjusted position. The housing joiner is then attached to the cable housing in front of the forward cable stop. Thereafter, the caliper brake is then reinstalled to its original installed position. The new, auxiliary, brake wire is extended from the housing joiner to interconnect with the main brake wire at a point forward of a second cable stop. System tension is reestablished by pulling the new cable toward the back of the bike, and connecting the two cables with a clamp. It is contemplated that installation of the auxiliary braking system may be easily and readily and rapidly achieved by a layperson. The system may likewise be readily and rapidly detached by reversing these installation steps, thereby restoring the bicycle to its original condition. More specifically, precise adjustment of the brake wire is critical to predictable braking characteristics, including braking tension and release clearances when unbraked, and maintaining this precision adjustment is of significance during attachment of the supplemental system of the present invention.\nIt should be noted and understood that with respect to the embodiments of the present invention disclosed herein, the materials, methods, apparatus and processes disclosed and suggested may be modified or substituted to achieve the desired protected structures without departing from the scope and spirit of the disclosed and claimed invention."}
{"text": "The invention relates to a catadioptric projection lens for imaging a pattern arranged in an object plane onto an image plane.\nProjection lenses of that type are employed on projection exposure systems, in particular wafer scanners or wafer steppers, used for fabricating semiconductor devices and other types of microdevices and serve to project patterns on photomasks or reticles, hereinafter referred to generically as “masks” or “reticles,” onto an object having a photosensitive coating with ultrahigh resolution on a reduced scale.\nIn order create even finer structures, it will be necessary to both increase the image-end numerical aperture (NA) of the projection lens to be involved and employ shorter wavelengths, preferably ultraviolet light with wavelengths less than about 260 nm.\nHowever, there are very few materials, in particular, synthetic quartz glass and crystalline fluorides, such as calcium fluoride, magnesium fluoride, barium fluoride, lithium fluoride, lithium-calcium-aluminum fluoride, lithium-strontium-aluminum fluoride, and similar, that are sufficiently transparent in that wavelength region available for fabricating the optical elements required. Since the Abbé numbers of those materials that are available lie rather close to one another, it is difficult to provide purely refractive systems that are sufficiently well color-corrected (corrected for chromatic aberrations). Although this problem might be solved by employing purely reflective systems, fabricating such mirror systems involves considerable expense and effort.\nIn view of the aforementioned problems, catadioptric systems that combine refracting and reflecting elements, i.e., in particular, lenses and mirrors, are primarily employed for configuring high-resolution projection lenses of the aforementioned type.\nWhenever imaging reflective surfaces are involved, it will be beneficial to employ beam-deflection devices whenever images free of obscurations and vignetting are to be achieved. Both systems having geometric beamsplitters having one or more fully reflecting deflecting mirrors and systems having physical beamsplitters are known. Additional plane mirrors may also be employed for folding the optical path. However, folding mirrors are usually employed only in order to allow meeting space-occupation requirements, in particular, in order to allow orienting the object and image planes parallel to one another. However, folding mirrors are not absolutely necessary from the optical-design standpoint.\nEmploying systems having a physical beamsplitter in the form of, for example, a beamsplitter cube (BSC), has the advantage that it allows configuring on-axis systems. Polarization-selective reflective surfaces that either reflect or transmit incident radiation, depending upon the orientation of its polarization axis, are employed in such cases. The disadvantage of employing such systems is that hardly any suitable transparent materials are available in the desired, large volumes. Moreover, fabricating optically active beamsplitter coatings situated within beamsplitter cubes may prove extremely difficult, particularly in cases involving large angles of incidence on the reflective surfaces involved and/or radiation that is incident over a broad range of angles of incidence.\nSome of the disadvantages of systems equipped with beamsplitter cubes may be avoided by employing systems having one or more deflecting mirrors in their beam-deflection devices. However, such systems have the inherent disadvantage that they are necessarily off-axis systems, i.e., systems with off-axis object fields.\nA catadioptric reduction lens of that type is disclosed in European Patent EP 0 989 434, which corresponds to U.S. Ser. No. 09/364,382. Such lenses have a catadioptric first lens section having a concave mirror and a beam-deflection device arranged between their object plane and their image plane and a dioptric second lens section that is arranged following that first lens section. Their beam-deflection device, which is configured in the form of a reflecting prism, has a first reflective surface for deflecting radiation coming from their object plane to the concave mirror and a second reflective surface for deflecting radiation reflected by the concave mirror to their second lens section, which contains exclusively refractive elements. A positive lens whose refractive power is dimensioned such that the concave mirror lies in the vicinity of the pupil is arranged between the object plane and the first reflective surface. The catadioptric first lens section creates a real intermediate image that is situated a short distance behind the second reflective surface and far in front of the first lens of the second lens section. This intermediate image is thus readily accessible, which may be taken advantage of for the purpose of, e.g., installing an illuminated-field stop. Large maximum angles of incidence, particularly at the first reflective surface, impose stringent demands on the reflective coatings used on the mirrors in order that they will have largely uniform reflectances for all incident radiation.\nAnother reduction lens that has a beam-deflection device having a deflecting mirror is disclosed in U.S. Pat. No. 5,969,882, which corresponds to European Patent EP A 0 869 383. In the case of this particular system, the deflecting mirror is arranged such that light coming from the lens' object plane initially falls on the concave mirror of its first lens section, which reflects it to the beam-deflection device's deflecting mirror, which reflects it to another reflective surface that deflects it toward the lens of the entirely dioptric second lens section. The elements of the first lens section that are utilized for creating the intermediate image are configured such that the intermediate image lies close to the beam-deflection device's deflecting mirror. The second lens section refocuses the intermediate image onto the image plane, which may be oriented parallel to the object plane, thanks to the reflective surface that follows the intermediate image in the optical train.\nU.S. Pat. No. 6,157,498 depicts a similar configuration whose intermediate image lies on, or near, the reflective surface of its beam-deflection device. Several lenses of its second lens section are arranged between the latter and a deflecting mirror located in its second lens section. An aspheric surface is also arranged in the immediate vicinity of, or near to, its intermediate image, and is included exclusively for the purpose of correcting for distortions, without affecting other imaging errors.\nA projection lens having a reducing catadioptric lens section and an intermediate image situated in the vicinity of the deflecting mirror of a beam-deflection device is depicted in German Patent DE 197 26 058.\nU.S. Pat. No. 5,999,333 depicts another catadioptric reduction lens having deflecting mirrors for which light coming from its object plane, after transiting a lens group having a positive refractive power, initially strikes a concave mirror that reflects it to the beam-deflection device's sole reflective surface. The intermediate image created by its catadioptric lens section lies close to this reflective surface, which reflects light incident thereon to a dioptric second lens section that images the intermediate image onto its image plane. Both its catadioptric lens section and its dioptric lens section create reduced images.\nA similarly configured lens for which the intermediate image created by its catadioptric lens section also lies close to its beam-deflecting device's sole deflecting mirror is depicted in Japanese patent application JP-A-10010429. The surface of the lens of the dioptric lens section that follows that lies closest to the deflecting mirror is aspherical in order that it may make a particularly effective contribution to correcting for distortion.\nOther lenses having off-axis object fields, a geometric beamsplitter, and a single concave mirror and single intermediate image followed by a dioptric lens section are known from U.S. Pat. No. 5,052,783 A, U.S. Pat. No. 5,691,802, and European patent application EP 1 079 253 A.\nSystems whose intermediate image lies near, or on, a reflective surface may be compactly designed. They also allow keeping corrections for the field curvatures of these off-axis illuminated systems small, which simplifies correcting for distortion. However, intermediate images that fall on a reflective surface can be problematical, since, in that case, flaws on the reflective surface will be sharply imaged onto the image plane, and since extremely high irradiance levels may occur at the reflective surface.\nCatadioptric systems that have beamsplitters generally have a group of lenses that is transited twice, i.e., are transited by light on its way from their object field to their concave mirror and its way from their concave mirror to their image field. In U.S. Pat. No. 5,691,802, it was proposed that this lens group have positive refractive power, which supposedly would allow employing a concave mirror with a smaller diameter. A system having a single positive lens that is transited twice situated near a deflecting mirror of its beamsplitter is disclosed in U.S. Pat. No. 6,157,498.\nThe high prices of the materials involved and limited availability of crystalline calcium fluoride in sizes large enough for fabricating large lenses represent problems, particularly in the field of microlithography at 157 nm for very large numerical apertures, NA, of, for example, NA=0.80 and larger. Measures that will allow reducing the number and sizes of lenses employed and simultaneously contribute to maintaining, or even improving, imaging fidelity are thus desired."}
{"text": "1. Field of the Invention\nThe present invention relates to a disc handling apparatus such as a recording and/or reproducing apparatus in which, for example, a disc shaped rotary record bearing medium is used.\n2. Description of the Prior Art\nIn the case of a recording and/or reproducing apparatus for recording and/or reproducing information with use of a record bearing medium such as a rotary magnetic disc, it is important to mount the record bearing medium on a rotation driving shaft correctly and surely in order to avoid eccentricity or the like. If the record bearing medium is not correctly mounted on the rotation driving shaft, the contact of a recording and/or reproducing head with the record bearing medium is not perfect (for example, a spacing loss increases). Or, in the case of eccentricity, a deviation occurs between the head and a track, and the recording or reproducing property becomes inferior, which is problematic.\nIn order to avoid the above difficulty, a method in which a pressing member is used to forcedly press the record bearing medium against a conical part which is provided on a record bearing medium rotating shaft has been proposed. However, according to this method, the record bearing medium is always against the conical part of the driving shaft and the pressing member rotates together with the record bearing medium, so that the moment of inertia of the rotating portion increases and also the size of a motor for driving the rotating shaft increases, which is disadvantageous for a compact apparatus.\nIn order to avoid the above disadvantage, as described later, an apparatus has been proposed in which: a plastic center hub having an engaging hole for engaging with a rotating driving shaft is provided at a central portion of a flexible magnetic disc, and a spring portion for side pressure is formed on a part of the hole, so as to obtain side pressure between the driving shaft and the center hub; and at a lower part of the center hub a magnetizable plate such as one made of iron is secured, while at a flange portion of the driving shaft for determining a height of a disc relative to a recording and/or reproducing head a permanent magnet is disposed, so as to effect chucking of the disc by the magnetic attracting force between the magnetizable plate and the permanent magnet.\nHowever, even in the case of the above method, for example, when the side pressure of the spring portion formed in the engaging hole of the center hub is low, eccentricity of the disc is apt to occur. Conversely, when the side pressure is high, chucking by the attracting force between the magnetic plate of the center hub and the permanent magnet at the side of the rotating driving shaft becomes uncertain. In order to obtain a sure chucking, the permanent magnet has to be made large, which means that the load for the motor is increased, and therefore the motor as well, has to be made large, which is disadvantageous."}
{"text": "Heterogeneous photocatalysis is the general term that describes the technical approach. Mills, A., Le Hunte, S., “An Overview of Semiconductor Photocatalysis,” J. PhotoChem. & PhotoBio. A: Chemistry 108 (1997) 1-35; and Hoffman, M. R., Martin, S. T., Choi, W., Bahnemann, D. W., “Environmental Applications of Semiconductor Photocatalysis,” Chem Rev 1995, 95, 69-96. Of particular importance is the formation of OH., the hydroxyl radical. The hydroxyl radical is an extremely potent oxidizing agent (redox potential of +2.8 eV vs. SHE (Standard Hydrogen Electrode)) that is capable of oxidizing almost all organic compounds. By comparison, the redox potentials for the more conventional oxidants chlorine and ozone are +1.36 and +2.07 eV, respectively. Hydroxyl radicals also kill and breakdown microorganisms and endotoxins.\nContaminants in fluid streams, such as organic compounds, nitrogen and sulfur oxides, acid gasses, dissolved inorganic solids, and microorganisms are converted by the oxidizing and reducing potential of the activated semiconductor. The conversion may take the form of, but is not limited to, the oxidation of organic compounds, degradation of microorganisms, or reduction of dissolved ionic species. The products of the conversion are ideally less harmful or more easily removed from the fluid stream than the parent compounds.\nSemiconductor photocatalysts that have been demonstrated for the destruction of organic contaminants in fluid media include but are not limited to: TiO2, ZrO2, ZnO, CaTiO3, SnO2, MoO3, Fe2O3, and WO3. TiO2 is the most widely investigated because it is chemically stable, has a suitable bandgap structure for UV/Visible photoactivation, and is relatively inexpensive.\nTiO2 exists as three principal crystalline forms: rutile, brookite, and anatase. The rutile form of TiO2 is widely used as a pigment and can be found in almost anything white, e.g., paint, paper, textiles, inks, plastics and cosmetics. Anatase, the low temperature form, is the most photoactive form. One method for making a brookite-containing TiO2 photocatalyst is described in U.S. Pat. No. 6,337,301, the disclosure of which is incorporated by reference in its entirety.\nThe inclusion of co-catalysts (e.g., platinum, palladium, silver and/or their oxides and sulfides) with titanium dioxide can increase the photocatalytic activity. A variety of methods improve the quantum efficiency of TiO2 by adding various metals to increase the minority carrier diffusion length, Augustynski, J.; Hinden, J. Stalder, C.; J. Electrochem. Soc. 1977, 124, 1063, or achieve efficient charge separation to increase carrier lifetimes. Vogel, R., Hoyer, P., Weller, H., “Quantum-Sized PbS, CdS, Ag2S, Sb2S3 and Bi2S3 Particles as Sensitizers for Various Nanoporous Wide-Bandgap Semiconductors,” J. Phys. Chem. 1994, 98, 3181-3188.\nThe inventors have found that the performance of a photocatalytic fluid purification device is primarily a function of the activity of the semiconductor photocatalyst, the specific surface area of the irradiated semiconductor surface, the mass transfer characteristics of the device, the pressure drop across the device, and light distribution and absolute light intensity at the semiconductor surfaces. All of these factors should be considered concurrently in the design of the photoreactor. To take advantage of the activity of the semiconductor surface, conditions are preferably controlled such that the rate of transport of contaminants from the fluid bulk to the semiconductor surface is not a limiting factor. Because the rate of contaminant conversion increases with increased surface concentration of the contaminant, the concentration gradient between the fluid bulk and the fluid at the semiconductor surface should be preferably minimized. The mass transfer characteristics of the device increase with the velocity of the fluid through the device and with closer proximity of the bulk fluid to the semiconductor surfaces. Preferred mass transfer characteristics are achieved by altering pore diameter, packing diameter, or otherwise physically ensuring small characteristic lengths. Forcing the fluid to flow in a tortuous path past the semiconductor surfaces also inherently improves mass transfer.\nAccording to the invention, the absolute conversion rate is a function of the surface concentration of the organic compound, the light intensity, and the irradiated surface area. The inventors believe that the reaction order in light intensity is <1, and that it is half order based on theoretical consideration of the accepted reaction mechanism and observed experimentally. Gerischer, H., Electrochimica Acta 38, 9 (1993) and Turchi, C. S., Ollis, D. F., Journal of Catalysis 119, 483 (1989). Based on the inventors' understanding that the reaction order for light intensity is less than unity, it is desirable to spread the light from a given source over as much surface area as possible. Therefore, according to the principles of the invention, an efficient device is one that provides intimate contact between the fluid and the semiconductor surfaces while spreading activating light over as large a surface area as is practical. Note that if the reaction order were, in fact, unity in light intensity, increasing the specific irradiated surface area, for a given light source, would not improve performance.\nSince until recently known semiconductor surfaces have been opaque to UV light; the spreading of light has typically been done by using geometry or waveguides. U.S. Pat. No. 5,516,492, to Dong, uses a plurality of curved plates while U.S. Pat. No. 6,063,343, to Say, uses a series of close packed plates orientated normal to the lamp axis to provide high specific irradiated surface area. U.S. Pat. No. 6,285,816, to Anderson, uses a waveguide to distribute the light. The previously mentioned methods rely on geometry or optics to distribute light to a semiconductor film and do not rely on a semitransparent substrate/semiconductor film pairing for radiation distribution.\nMuch of the early research on semiconductor photocatalysis concerned methods using opaque titanium dioxide (TiO2) slurries or TiO2 wash coatings onto or inside a glass tube and the photodegradation of organic compounds and their intermediates in water. These methods of using TiO2 have limitations for commercial applications. For example, TiO2 slurry has the serious limitation of the removal of the TiO2 particles from the purified water. While wash coating TiO2 onto glass avoids the removal limitations of the slurry approach, it has its own problems in that insufficient surface area is presented for effective destruction of organics within a reasonable time period. Additionally, the wash coat is not firmly attached to the glass resulting in a loss of TiO2 to the water stream and a concomitant reduction in photocatalytic activity.\nKraeutler and Bard made a photocatalytic reactor with a water slurry of suspended TiO2 powder, in the anatase crystalline form, and studied the decomposition of saturated carboxylic acid. J. ACS 100 (1978) 5985-5992. Other researchers used UV-illuminated slurries of TiO2 to study the photocatalyzed degradation of organic pollutants in water.\nMathews created a thin film reactor by wash coating TiO2, (Degussa P25™), particles to the inside of a 7 millimeter long borosilicate glass tube wound into a 65-turn spiral. He monitored the destruction of salicylic acid, phenol, 2-chlorophenol, 4-chlorophenol, benzoic acid, 2-naphthol, naphthalene, and florescin in water. J. Physical Chemistry 91 (1987) 3328-3333. U.S. Pat. No. 5,766,455, to Berman, wash coated an opaque coating onto a glass fiber mesh. Neither of these devices uses a transparent substrate/semiconductor pairing.\nU.S. Pat. No. 4,892,712, to Robertson et al., discloses the attachment by the sol-gel process of anatase TiO2 to a fiberglass mesh substrate. This mesh was wrapped around a light source contained within a quartz glass cylinder and emitting UV radiation in a wavelength range of 340 to 350 nanometers (nm). Unlike the present invention, Robertson's mesh is not rigid, lacks permanent bonding of the semiconductor to the mesh, and does not specify the transparency of the substrate semiconductor pairing.\nProfessor I. R. Bellobono prepared photocatalytic membranes immobilizing 23% of Titanium Dioxide (Degussa P-25). Controlled amounts of appropriate monomers and polymers, containing the semiconductor to be immobilized and photoinitiated by a proprietary photocatalytic system was photografted onto a non-woven polyester tissue. The final porosity of the photosynthesized membrane was regulated at 2.5-4.0 microns. He trade named this membrane “Photoperm”™. “Effective Membrane Processes. New Perspectives,” R. Paterson, ed., BHR, Mech. Eng. Publ., London (1993), 257-274. The process was patented in Italy in 1995, Italian Pat. No. IT1252586. Unlike the present invention, Bellobono's apparatus is not inert, not durable, and would display excess pressure drop due to low fluid permeability.\nCittenden, et al. discloses a method and apparatus for destroying organic compounds in fluids. See The 1995 American Society of Mechanical Engineers (ASME) International Solar Energy Conference, Maui, Hi., USA. An opaque coating of TiO2 was attached by wash coating to a 35×60-mesh silica gel substrate. The substrate was placed within a plastic tube that allowed the penetration of UV light. Organic pollutants in a water stream passed axially through the tube. Unlike the present invention with a semitransparent semiconductor coating, Cittenden's invention is not durable, the photocatalytic coating is not semitransparent, and has very limited fluid permeability.\nAnother method to make ceramic titanium membranes uses a refined sol-gel process. J. Membrane Science 30 (1988) 243-258, and U.S. Pat. No. 5,006,248, to Anderson. These membranes are porous and transparent to UV illumination. They are made from a titanium alkoxide and then fired to form the anatase crystalline structure. Unlike the present invention, Anderson's invention has very limited fluid permeability.\nThus, while attempts have been made in the prior art to enhance quantum yields by increasing semiconductor specific irradiated surface area and improving UV light penetration, serious limitations remain to the commercial development of an efficient, durable, photocatalytic purification apparatus with acceptable pressure drop characteristics and methods for its use."}
{"text": "The present invention relates to mastitis, an extremely important veterinary disease, and to the preparation and administration of egg antibody products which may be used to neutralize a wide variety of systemic pathogens which cause or exacerbate this disease.\nThe term mastitis refers to inflammation of the mammary gland, regardless of cause. It is characterized by physical, chemical, and usually bacteriological changes in the milk and by pathological changes in the glandular tissue. The most important changes in the milk include discoloration, the presence of clots and the presence of large numbers of leukocytes. Although there is swelling, heat, pain, and induration in the mammary gland in many cases, a large proportion of mastitic glands are not readily detectable by manual palpation or by visual examination of the milk using a strip cup. Because of the very large numbers of such \"subclinical\" cases, the diagnosis of mastitis has come to depend largely on indirect tests which depend, in turn, on the leukocyte content of the milk. Thus, it is practicable to define mastitis as a disease characterized by the presence of a significantly increased leukocyte content in milk from affected glands. This leukocyte content is commonly referred to as the \"somatic cell count,\" or SCC.\nMany infective agents, primarily bacterial but occasionally fungal or protozoal, have been implicated as causes of mastitis. The most common are Streptococus agalactiae and Staphylococcus aureus, with Escherichia coli becoming a significant cause in housed or confined cattle, principally in the northern hemisphere. Other agents include Streptococcus uberis, Streptococcus dysgalactiae, Corynebacterium (Actinomyces) pyogenes, Klebsiella spp., Mycobacterium bovis, Pseudomonas pyocyaneus, Serratia marcescens, Mycoplasma bovis, Nocardia asteroides, Proteus spp., and Chlamydia psittaci.\nInfection of each mammary gland occurs via the teat canal, with the infection originating from two main sources, the infected udder and the environment. In dairy cattle and milking goats, the important infections are those which persist readily in the udder, especially Str. agalactiae and Staph. aureus. The contamination of milkers' hands, washcloths, and milking machine cups by milk from infected quarters may quickly lead to the spread of infection to the teats of other animals. These are commonly referred to as infectious or contagious mastitides.\nBacteria which are normal inhabitants of the environment, such as E. coli and Ps. pyocyaneus, cause mastitis less frequently but, when they do, the disease is much more resistant to hygienic control measures. Large, zero-grazed herds in feedlots are likely to encounter more hygiene problems than conventionally housed herds. The problems relate mainly to fouling of the udder caused by inadequate or improper bedding in the larger units. Thus, these are known as environmental mastitides.\nThe frequency of occurrence of mastitis is dependent upon the ability of bacterial agent to set up infection in the mammary tissue. The diferences between bacteria in their ability to set up a mastitic state is dependent on at least two important groups of factors: bacterial characteristics and transmission mechanisms. Bacterial characteristics include the ability of the organism to survive in the cow's immediate environment (i.e., its resistance to environmental influences, including cleaning and disinfection procedures), its ability to colonize the teat duct, its ability to adhere to mammary epithelium and set up a mastitis reaction, and its resistance to antibiotic therapy. Transmission mechanisms are a function of the extent of the infection in the environment (including infected quarters), the efficiency of milking personnel and machinery, milking parlor hygiene, and the susceptibility of the cow. An individual cow's susceptibility, in turn, is a function of the stage of lactation (early lactation, especially the first two months, representing highest susceptibility), the age of the cow (older cows, having had more than four lactations, being more susceptible), the level of inherited resistance (possibly related to teat shape and anatomy of the teat canal), lesions on the teat skin (especially the orifice), and the immunological status of each mammary gland (including prior infections, especially with Staph. aureus).\nIn most countries, surveys of the incidence of mastitis, irrespective of cause, show comparable figures of about 40% morbidity amongst dairy cows and a quarter infection rate of about 25%. The incidence is similar in goats. Surveys of the prevalence of the various infections in cattle show remarkable similarity in different countries. The predominant position of Str. agatactiae as a cause of bovine mastitis has been usurped by Staph. aureus, especially in areas where the treatment of mastitis with penicillin has been practiced intensively and where machine milking has replaced hand milking. In such areas a relative incidence of Str. agalactiae, other streptococci, and Staph. aureus of 1:1:2 is a common finding.\nAlthough the prevalence of Staph. aureus has been significantly curbed by modem control programs based on teat dipping and dry period treatment, it remains the preeminent cause of subclinical mastitis. Mastitis is said to be subclinical when there is evidence of inflammation, e.g., a high somatic cell count in the milk, without any visible abnormality of the milk or udder. Chronic cases of mastitis, particularly those caused by Staph. aureus, are often refractory to treatment while the cow is lactating. For staphylococcal infections in general, a cure rate of 30% is about the best that can be expected. An old adage holds that \"once a Staph. cow, always a Staph. cow.\" Thus, many dairymen choose to cut their losses and simply cull Staph.-positive cows rather than to attempt treatment.\nAlthough mastitis occurs sporadically in all species, it assumes major economic importance only in dairy cattle. In terms of economic loss it is undoubtedly the most important disease with which the dairy industry must contend. In 1992, mastitis infections cost members of the American Dairy Association over $200,000,000. High somatic cell counts in cows result in the loss of approximately 200,000 pounds of milk per cow per year at a financial cost of approximately $288 per cow. Only 10% of the approximately 9,750,000 dairy cows in the United States attain the highest ranking on somatic cell count (below 100,000). A dairyman receives an additional 50 cents per 100-wt. of milk for the highest ranking, i.e., the lowest SCC range. Thus, the insidious presence of subclinical mastitis in a herd can have a devastating effect on milk profits, even in the absence of any external manifestations of disease.\nMost estimates show that on the average an affected quarter suffers a 30% reduction in productivity and an affected cow is estimated to lose 15% of its production. Total economic losses caused by mastitis comprise the value of milk production lost, the value of cows lost by premature culling, the value of milk discarded or downgraded, and the cost of treatment and veterinary expenses.\nOrdinarily, mastitis is treated either systemically by parenteral injection or locally by intramammary infusion. Parenteral treatment is usually reserved for cases of mastitis in which there is a marked systemic reaction, to control or prevent the development of a septicemia or bacteremia and to assist in the treatment of the infection in the gland. The systemic reaction can usually be brought under control by standard doses of antibiotics or sulfonamides but complete sterilization of the affected quarters is seldom achieved because of the relatively poor diffusion of the antibiotic from the bloodstream into the milk. Parenteral treatment is also used when the gland is badly swollen and intramammary antibiotic is unlikely to diffuse properly. To produce therapeutic levels of antibiotic in the mammary gland by parenteral treatment, it is necessary to use higher-than-normal dose rates, which can alter milk-withholding times, sometimes with disastrous consequences, as discussed infra.\nBecause of convenience and efficiency, udder infusions have been the preferred method of treatment. Disposable tubes containing suitable drugs in either an aqueous or ointment base are inserted into the teat canal, providing diffusion of drug into the glandular tissue. Strict hygiene is necessary during treatment to avoid the introduction of bacteria, yeasts and fungi into the treated quarters. Diffusion of infused drugs is often impeded by the blockage of lactiferous ducts and alveoli with inflammatory debris. Drug-withholding times are just as applicable to udder infusions as parenterally administered drugs, since these infusions rapidly undergo systemic absorption from the highly vascular udder.\nChoice of a suitable drug for the treatment of mastitis is often extremely difficult. In vitro laboratory testing of bacterial sensitivity is not necessarily a justifiable basis for selecting the antibacterial agent to be used in individual cows, and the response to treatment in clinical cases is often unrelated to the results of in vitro sensitivity tests.\nIn the opinion of the authors of one of the leading texts on large animal veterinary medicine, very little good research and development work has been performed in the field of mastitis prophylaxis and treatment, and most of the products available have been developed with very little scientific support. Attractiveness to the consumer has been placed above efficiency, in the opinion of these authors, resulting in \"intramammary preparations containing a battery of antibacterial agents, one or more of a long list of generally unhelpful adjuvants, and inappropriate recommendations about their use, especially with regard to frequency of administration and length of time after treatment for which milk must be withheld from sale.\" Further, \"[t]he indiscriminate and improper use of antibiotics for the treatment of clinical mastitis in lactating cows which has resulted from inaccurate promotion and uninformed use by farmers has meant that the control of mastitis has really received little assistance. Drug-resistant organisms, especially Staph. aureus, have been encouraged, and the chances of human consumers of dairy products being exposed to antibiotic residues has been increased.\" (Blood and Radostits, Veterinary Medicine: A Textbook of the Diseases of Cattle, Sheep, Pigs, Goats and Horses, 7th ed., 1989, p. 512).\nThe effect of antibiotics in milk on the manufacture of dairy products and the development of sensitivity syndromes in human beings is of the utmost consequence. In most countries, the maximum intramammary dose of antibiotics is limited by legislation and the presence of detectable quantities of antibiotics in milk constitutes adulteration. Strict legislation and mandatory milk-testing in the United States have greatly decreased such risk to consumers, but the penalties for milk contamination as a result of inadequate milk-withholding times are enormous and can result in financial disaster for the dairyman and/or veterinarian found to be responsible for such a violation. Since the rate of drug excretion varies among different animals and in the same animal at different points in the lactation period, withholding times specified on drug labels are not always accurate and can lead to inadvertent violation of residue limits. As a result, some dairymen will elect to cull any but the most valuable cows rather than to run the risk of contaminated milk. Even aside from such risk, substantial loss is involved in the discarding of milk during specified withholding times, which can range anywhere from 72 hours for an udder infusion in a lactating cow to 10 days for parenteral antibiotics in long-acting bases.\nFor the foregoing reasons, there exists a felt need for a mastitis treatment which is not only efficacious, easily administered, and cost-effective, but which is not hampered by the uncertainty and risk of a milk-withholding period. Such a treatment would undoubtedly be enthusiastically received by the dairy industry worldwide."}
{"text": "This invention relates to a method and an apparatus for cutting off the beads of seam-welded pipes and especially of hot-shaped, seam-welded pipes More particularly, it relates to a method and apparatus for cutting off the beads of seam-welded pipes using a freely-rotatable circular cutting tool.\nIn the manufacture of seam-welded pipes and particularly hot-shaped seam-welded pipes, an endless strip, referred to as a skelp, is continuously passed through a heating furnace and heated to a temperature of at least 850.degree. C., after which it is passed through a series of forming rolls which bend the outer edges of the skelp and form it into the shape of a tube. The resulting tubular object is referred to as an open pipe. The confronting outer edges of the open pipe are as of yet unwelded. The open pipe is then passed through a tubular work coil which is somewhat larger than the outer diameter of the open pipe, and the confronting edges of the open pipe are selectively heated to a welding temperature, and the confronting edges of the open pipe are then pressure welded by squeeze rolls which are located downstream of the work coil, thereby forming a seam-welded pipe. Due to the heating and pressure welding, the inner and outer surfaces of the welded seam swells to form beads on the inside and outside of the pipe. The pipe is then sent to a finishing apparatus which cuts off the beads. It is important to cut off the beads smoothly so that there is no difference in thickness between the portion of the pipe which is cut and the adjoining portions. However, as the removal of the inner bead must be performed inside the pipe, it is generally difficult to cut the beads off smoothly. Normally, a slender mandrel is inserted from an open space upstream of where the forming rolls form the skelp into a tube to the downstream side of the squeeze rolls, and a circular cutting tool is mounted on the tip of the mandrel. Accordingly, the mandrel passes through the entire length of the open pipe and the cutting tool is secured to the mandrel downstream of the squeeze rolls.\nIn order to selectively heat the edge portions of the open pipe which are to be welded, a high-frequency current is passed therethrough by the work coil. Due to the skin effect, currents are induced on the circumferential surface of the open pipe, and heating is accomplished by the concentration of the current in the portions to be welded. However, if the induced current leaks to outside the portions to be welded, the welding efficiency decreases. Therefore, the leak current should be as small as possible, and in actual practice, an impeder having a ferrite core is disposed inside the open pipe so as to confront the work coil. Generally, since a mandrel for cutting off the beads on the inside of a seam-welded pipe passes through the center of the open pipe, the impeder is mounted on the mandrel opposite the work coil.\nSeam welding can be formed in a hot or cold state. When cutting off the beads on the inside of a cold-shaped, seam-welded pipe, in order to prevent excess or inadequate cutting, a full-type cutting tool having the same curvature as the inside of the pipe is used. Even when the seam-welded pipe is itself in a cold state, as cutting is performed immediately after welding when the bead is still hot, the cutting tool normally has to be replaced after 1.5-2.5 hours of use. The production line must be stopped for a long time when the cutting tool is replaced. The time lost for replacement of the cutting tool has an adverse effect on productivity, since the forming and welding processes for seam-welded pipes and the bead cutting operation are all performed continuously, so that when the cutting tool is replaced, all three of these operations must be halted.\nWhen manufacturing hot-shaped seam-welded pipes, if the same cutting method is used as for cold-shaped seam-welded pipes, decreases in productivity due to cutting tool replacement are even greater. When replacing the cutting tool, it is necessary to stop the production line just as when manufacturing cold-shaped seam-welded pipes. While the production line is stopped, it is also necessary to lower the temperature of the furnace in which the skelp is heated to the rolling temperature so that the skelp will not melt through and sever. After the cutting tool is replaced and the line is restarted, the furnace temperature is again raised, but skelp must be passed through the furnace for some length of time before the skelp temperature stabilizes, and the length of skelp which passes through the furnace until the skelp temperature stabilizes must be discarded. Furthermore, as the cutting tool operates in a high-temperature environment, the life span of the cutting tool and the length of time between replacements of the cutting tool are much shorter than in the manufacture of cold-shaped seam-welded pipes. Accordingly, if the cutting tool must be frequently replaced, not only are there large time losses, but the energy consumption rate rises and there is a large amount of material loss in the form of scrap.\nMoreover, when the inside beads of hot-shaped seam-welded pipes are cut off on a production line, just as with cold-shaped seam-welded pipes, a mandrel is inserted into a pipe from a position upstream of where the skelp is formed into an open pipe and the inside bead is cut off by a cutting tool mounted on the end of the mandrel. In hot shaping, the temperature of the skelp is at least 850.degree. C., so even without the heat generated by welding, the mandrel would be heated to around 700.degree. C. by the heat from the open pipe. Therefore, it is generally necessary to provide some form of heat protection for a mandrel used to cut off the inside beads of hot-shaped seam-welded pipes.\nFor these reasons, in a production line for seam-welded pipes and particularly hot-shaped seam-welded pipes, there is a great need to increase the life span of the cutting tool used for cutting off inside beads and to cool the mandrel on which the cutting tool is mounted.\nThere have been various proposals for using a circular cutting tool for removing the inside beads of seam-welded pipes. Japanese Published Unexamined patent application No. 48-56547 (1973) discloses the use of a circular cutting tool for removing the inner beads of cold-shaped seam-welded pipes. The cutting tool has a skew angle of 22.5.degree., although the diameter and the tilt angle of the cutting tool are not specified. No mention is made of the life span of the cutting tool.\nJapanese Published Unexamined patent applications Nos. 58-143908 (1983) and 59-144586 (1984) also disclose the use of circular cutting tools for removing inside beads of hot-shaped seam-welded pipes. By using a cutting tool with a diameter of 0.6-0.8 times the diameter of the pipe, a skew angle of 10.degree.-45.degree., and a tilt angle of 5.degree.-15.degree., a long life span of 2.5 hours (application No. 58-143908) or 5 hours (application No. 59-144586) can be achieved. For cutting in a high-temperature environment, these life spans are extremely long.\nIn above-mentioned Japanese Published Unexamined patent applications Nos. 58-143908 and 59-144586, in order to protect the mandrel from heat, the mandrel is formed with a hollow center, and cooling water is passed through the entire length of the mandrel from the base end and then continuing through the holder, the impeder, and the outer end, i.e., head portion on which the cutting tool is mounted. The cooling water is discharged into the seam-welded pipe through gaps between the outer end, i.e., head portion of the mandrel and a plurality of rollers which guide the mandrel along the inside beads of a pipe.\nAs for the cutting tool itself, in above-mentioned Japanese Published Unexamined patent application No. 48-56547, if the tilt angle is 5.degree.-15.degree., at a relatively slow line speed of at most 40 m/min, for example, satisfactory cutting can be carried out. However, when the line speed exceeds 100 m/min, the rotation of the cutting tool becomes unstable, and the cutting tool tends to remove too much or too little of the inside surface of a pipe being processed.\nIn Japanese Published Unexamined patent applications Nos. 58-143908 and 59-144586, the same tendency can be seen when the skew angle is small, but when the skew angle is large, the shape of the cutting tool as seen from the direction of advance of the pipe becomes elliptical, the difference between the curvature of the inside of the pipe and the cutting tool increases, and the cut surface takes on a U-shaped cross section, so that it becomes difficult to perform smooth cutting.\nAs the life span of cutting tools increases, the length of time which the mandrel must be inserted into an open pipe necessarily increases, and problems occur which were not experienced when cutting tools had a shorter life span.\nThe holder of the mandrel must support the mandrel so that the opposite end of the mandrel is located downstream of the squeeze rolls, and it must support an impeder so that the impeder confronts a work coil. Therefore, the holder is generally long, and the length of time for which cooling water passes therethrough is also long. As it passes through the holder, the cooling water undergoes a great increase in temperature, and even if the cooling water does not boil, gases which are dissolved in the cooling water come out of solution and remain inside the holder. The longer the mandrel is inserted into an open pipe, the larger the amount of gas which will accumulate inside the holder, and as the gas does not conduct heat as well as the cooling water, a difference in the rate of cooling will develop between the top side of the holder, where the gas accumulates, and the bottom side of the holder, which is in contact only with cooling water. This difference in cooling rates causes the holder to bend. Even if the amount of curvature is minute, as the holder itself is long, the misalignment between both of its ends and the center of the holder is large. Not only does this curvature affect the depth of cutting of the cutting tool, in the worst case, it becomes impossible to detach the rest of the mandrel from the pipe.\nThe impeder of the mandrel contains a ferrite core which is self-heated by magnetic induction caused by the action of the work coil. Accordingly, the greater the ability to cool the impeder the better. However, the gas which is generated in the holder is carried downstream to the impeder, and due to the gas, the impeder may be inadequately cooled. If the thermal load on the impeder is too large, not only will the ferrite core lose its magnetism, but the FRP cover which surrounds the ferrite core can be damaged, possibly resulting in the breakage of the impeder."}
{"text": "Known hydroelectric generating equipment systems have mechanical drawbacks which substantially limit their economic utilization - particularly at low head sites. Particularly, all known types require a direct mechanical connection between the turbine and the generator. Known low-head axial turbine equipment are of three types; (1) The rim-generator type--in which the generator rotor is located on the periphery of the turbine runner; (2) the tube-type--in which the generator is located outside the water passage and connected to the turbine through an extended drive shaft; and (3) the bulb turbine type--in which both the runner and the generator are directly connected and enclosed within the water passage.\nAll three types require expensive complicated concrete Civil Works to encase and structurally support the turbine and generator in a common housing. There are also problems in matching up turbine speeds with generator speed requirements, accessibility for maintenance purposes to the important components and the inability to install most known equipment into most already existing dams or structures.\nThe following described invention proposes to solve all of the above problems and more, by allowing a physical separation of the turbine and the generator and accomplishing several other required features (such as synchronous generator speed control) through the unique application of certain fluid drive components."}
{"text": "The present invention relates to wellbore tools and more particularly to an apparatus for retrieval of metal objects, such as cuttings and other foreign objects that accumulate in the process of perforating or milling over bridge plugs and other down hole obstructions from a wellbore.\nVarious types of bridge plugs are conventionally used in the oil and gas industry. These bridge plugs are installed in the annulus and are often covered in cement. Removal of such plugs can sometimes pose a problem for the industry. A rotary bit drills the cement and plugs out, while some of the cuttings of the plugs are carried out to the surface by a liquid circulated down hole.\nSometimes, a production packer needs to be removed together with the metal pipe that it surrounds. In those cases, milling tools with gravity fed boot baskets are used for retrieving pieces of metal from the wellbore. After retrieval of the production packer, it may become necessary to run a conventional fishing magnet to retrieve additional junk and cuttings.\nA conventional fishing magnet is mounted inside a housing that is lowered into a wellbore. It is limited in the ability to retrieve cuttings in that its magnetization is restricted to the extreme bottom surface of the magnet. The fact that circulating fluids lift the cuttings away from the bottom surface of the magnet renders that conventional fishing magnet useless in this situation.\nOften times, a boot basket is used for collecting cuttings that did not attach themselves to the conventional magnet. A boot basket has small openings for catching these particles. Consequently, many large size pieces or very small pieces suspended in the fluid flow are not trapped in the basket and remain in the wellbore.\nThe present invention contemplates elimination of drawbacks associated with the prior art and provision of a wellbore apparatus for removal of metal objects, such as cuttings, and other foreign particles that provides for the use of sets of magnets spaced longitudinally along the tool body.\nIt is, therefore, an object of the present invention to provide an apparatus for retrieval of metal cuttings and other foreign objects from a wellbore.\nIt is another object of the present invention to provide a packer-milling tool with magnets set in a spaced relationship longitudinally along the tool body to increase xe2x80x9ccatchingxe2x80x9d capacity of the tool.\nThese and other objects of the invention are achieved through a provision of an apparatus for removal of metal cuttings and other foreign objects from a wellbore that uses at least two magnet assemblies vertically spaced from each other and mounted on the body of the tool. Each magnet assembly is comprised of a plurality of magnets encircling the body of the tool and covering a surface area greater than one half of the tool body.\nOne of the embodiments of the invention provides for arcuate in cross section magnet members that are detachably secured on the tool body. The second embodiment provides for trapezoidal in cross section magnet members that extend longitudinally in parallel relationship to a central axis of the tool body. The tool of the second embodiment further provides for magnet protectors mounted adjacent to each magnet member. Each magnet protector has a triangular cross section and has a slanted surface that is angled in the direction of rotation.\nThe space between each magnet and the next adjacent magnet protector forms a xe2x80x9ctrapxe2x80x9d that receives and retains metal particles pushed away by the leading surface of the adjacent magnet member. Consequently, the capability of the tool to retain cuttings for subsequent retrieval to the surface is significantly increased in comparison with conventional magnet tools."}
{"text": "Field\nThis invention relates generally to transceivers, and more specifically, to a reconfigurable multi-mode transceiver.\nBackground\nA transceiver in a communication system may transmit data using a transmitter and receive data using a receiver for two-way communication. The transceiver may also support carrier aggregation (CA), which is simultaneous operation on multiple carriers. Multi-mode transceivers including Long Term Evolution (LTE)/Frequency Division Duplex (FDD). LTE/Time Division Duplex (TDD). Circuit-Switched FallBack (CSFB), Simultaneous Voice and LTE (SVLTE), Simultaneous GSM and LTE (SGLTE), Dual SIM Dual Standby (DSDS), and other similarly-configured devices co-exist in several different combinations in different regions of the world. Further, LTE carrier aggregation with 2, 3, and 4 downlinks is already out or will soon be available. Thus, a device that can handle multiple modes and bands while providing seamless FDD and TDD interoperability is needed in order to meet the growing demands for wireless data. However, the support of both FDD and TDD operations complicates the local oscillator LO and synthesizer connections to the receive (RX) and transmit (TX) links in the multi-mode transceiver.\nTo address the issues connected with a need for multi-band, multi-mode devices, a two-chip configuration in which one transceiver (e.g., an FDD transceiver) resides on a first chip and another transceiver (e.g., a TDD transceiver supporting LTE-CA and other technologies) resides on a second chip can be used. Although this configuration may provide the needed features, it will occupy more chip area, consume more current, and support less number of combinations of CA than a single-chip solution. Accordingly, a need exists for reconfiguring and reusing the transceiver flexibly with minimum hardware to provide a single-chip solution for improving the FDD and TDD interoperability and better spectrum utilization."}
{"text": "The present invention relates generally to methods for treatment of pulmonary diseases and diseases associated with upper respiratory tract infection. Such diseases, including cystic fibrosis, emphysema, chronic bronchitis, sinusitis, and the common cold, have in common bronchial or sinus congestion, production of large amounts of sputum, and the possibility of secondary bacterial infection requiring antibiotic therapy.\nAcute otitis media is a bacterial or viral infection in the middle ear which is usually secondary to upper respiratory tract infections and is most common in children. Microorganisms may migrate from the nasopharynx to the middle ear over the surface of the eustachian tube's mucous membrane or by propagating in the lamina propria of the mucous membrane as a spreading cellulitis or thrombophlebitis. Pain and hearing loss are the most common presenting complaints although fever, nausea, vomiting and diarrhea may occur in young children. Therapy for acute otitis media includes analgesics, decongestants and antibiotics. In addition, topical vasoconstrictors may be administered into the nasal cavity to improve eustachian tube function. Further, systemic sympathomimetic amines such as ephedrine sulfate may also be administered.\nSerous otitis media (secretory otitis media) is an effusion in the middle ear resulting from incomplete resolution of acute otitis media or obstruction of the eustachian tube. Traditional therapy includes a trial of antibiotic therapy in case of bacterial infection. Such antibiotic therapy is effective in relieving eustachian tube obstruction due to bacterial infection and in sterilizing the middle ear. Systemic sympathomimetic amines may also improve eustachian tube function by their vasoconstrictive effects and antihistamines may relieve eustachian tube obstruction in allergic patients. Surgical therapies include myringotomy for aspiration of the fluid and for insertion of a tympanostomy tube which allows ventilation of the middle ear and ameliorates the eustachian tube obstruction. Alternatively, the middle ear may be temporarily ventilated with the Valsalva maneuver or politzeration.\nDespite the efficacy of these approaches there remains a desire to avoid surgical intervention in cases of otitis media. Moreover, there exists a growing concern that the widespread use of antibiotics for treatment of otitis media in children promotes the development of antibiotic resistant bacteria. Accordingly, there remains a desire in the art for improved treatment of conditions associated with upper respiratory infections and pulmonary disorders including otitis media."}
{"text": "This invention relates to a range detecting device which detects or measures a distance between an object and the device. Various types of range detecting devices for enabling automatic or semi-automatic focus adjustment of camera lenses are conventionally known, and they are generally classified as passive or active types. The passive type has merit in that the mechanism is simple, but it has a disadvantage in that range detection is difficult when the scene is dark.\nHowever, the active type of detecting devices projects the light from the device toward the object, and the reflected light is sensed, so that the disadvantage of the passive type devices is not present in this type. A camera having this active type of a range detecting devices is disclosed in U.S. Pat. No. 3,443,502.\nIn cameras using active devices such as the type shown in U.S. Pat. No. 3,443,502, a light source emits a beam which is moved angularly to sweep the field of the camera lens. The angular motion of the beam is linked to and calibrated with concomitant motion of the lens-focus-adjustment. A sensor is positioned to receive light reflected from the object, and, upon receiving this reflected light, the displacement of the focus-adjustment is arrested, thus positioning and holding the camera lens at the correct range adjustment.\nIn range detecting devices of the active type mentioned above when the light sensor fully receives the reflected light from the object and produces a signal of high intensity, the camera lens is at the correct range adjustment, and when the light sensor does not receive the reflected light and produces a signal of low intensity, the camera lens is not at the correct range adjustment. Therefore, for increasing the accuracy of the range detection, it is desirable that the sensor have a small displacement from the position of the beam which produces a low output signal to the position of the beam which produces a high output signal, namely, that the intensity of the sensor output signal vary definitively from a low state to a high state while the beam emitted from the light source is angularly moving to sweep the field of the camera lens.\nIt has been also desired to increase the maximum detectable distance that the active type range detecting device may detect between the object and the device."}
{"text": "1. Technical Field\nThe present invention relates to lighting devices.\n2. Description of Related Art\nWhen the lighting device is applied to the car lamp, the illuminative range and intensity of the lighting device should conform to the traffic standard and is provides safety while driving at night or in dark places such as tunnels.\nIn general, the lighting device includes a lens, a light source and a gobo. The gobo is located between the lens and the light source and controls the shape of the emitted light. But, the light from the light source is focused at the focal point by the lens. It's difficult to adjust the illuminative pattern of the lighting device."}
{"text": "The present invention relates to a hydraulic brake system with anti-lock control. Such brake systems include a pedal-operated pressure generator, normally a tandem master cylinder actuated by a booster-aided brake pedal. The master cylinder is connected to wheel brakes via at least one brake line which includes an electro-magnetically operated, normally open two-way, two-position valve. For pressure relief during anti-lock control operation, each of the wheel brakes is also connected to a low-pressure accumulator or reservoir via a pressure relief line. In the pressure relief line, an outlet valve is inserted which is a normally closed, electro-magnetically operated two-way, two-position valve. To return the fluid released from the wheel brakes into the brake circuit, a pump is connected with its suction side to the low-pressure accumulator or the fluid reservoir respectively, and with its pressure side through a pressure line to the brake line.\nAfter assembling such a brake system, the entire circuitry has to be evacuated and filled with brake fluid. In order to enable a non-powered evacuation and filling of such a brake system, a check valve with a low opening pressure was placed to connect the pressure relief line to the pressure line of the pump. The check valve would open in the direction enabling a pressure flow from the relief line to the pressure line, thereby permitting the evacuation of the pressure relief line during evacuation of the pressure line. Vacuum will be kept in the pressure relief line after filling the system with brake fluid. This area will be filled with brake fluid at the end of testing the brake system during activation of the anti-lock control, thus forcing brake fluid into the pressure relief line while the normally closed valves are energized and the pump is running. While this constellation worked well for the purpose of evacuation and fill, however, it caused vacuum trapping in the return circuitry, which is the pressure relief line, after installation of the hydraulic circuitry including a valve block into a vehicle. The biggest restrictions in the brake line are the normally open electro-magnetically operated valves. When a driver releases the brake pedal, in brake systems according to the prior art, vacuum can be created between the normally open valves of the valve block and the tandem master cylinder since pressure energy between the wheel brakes and the normally open valves changes to kinetic energy at the location of the normally open valves. When this vacuum creates a sufficient pressure difference between the brake line and the pressure relief line, the bleed valve opens and traps vacuum in the pressure relief line. After a while, the vacuum will be replaced by air due to small leakage in the pressure relief line of the hydraulic circuitry, and air accumulates after a certain number of brake operations. The result is a long brake pedal stroke or a reduced braking force during electronic control of the brake operation such as anti-lock control, traction control, or driving stability control during cornering. By raising the opening pressure of the bleed valve, not as much vacuum can be trapped, however, this would result in an unsatisfactory evacuation operation.\nThe object of the current invention therefore is to create a brake system which enables a non-powered evacuation and filling of the circuitry on the one hand and which prevents the trapping of vacuum on the other hand.\nAnother object of the invention is to create a low-cost solution to the above-mentioned problem without adding excessive additional equipment.\nThese objectives are achieved by a brake system of the above-mentioned kind where a bypass line connects the pressure relief line with the wheel brake. This bypass line includes a bleed valve opening from the pressure relief line to the wheel brake. A possible point for connecting the bypass line including the bleed valve to the brake line is any point in direct connection with the brake line. That means if the bypass line is connected to the pressure line of the pump without any other hydraulic element in-between, that would be equivalent to connecting the bypass line directly to the brake line. This solution provides a reliable vacuum compensation during pedal releases while having a low opening pressure for the bleed valve in order to evacuate the pressure relief line properly."}
{"text": "With an increased interest in cell-specific drug testing, diagnosis, and other assays, systems that allow for individual cell isolation, identification, and retrieval are becoming more desirable within the field of cellular analysis. Furthermore, with the onset of personalized medicine, low-cost, high fidelity cellular sorting systems are becoming highly desirable. However, preexisting cell capture systems suffer from various shortcomings that prevent widespread adoption for cell-specific testing. For example, flow cytometry requires that the cell be simultaneously identified and sorted, and limits cell observation to a single instance. Flow cytometry fails to allow for multiple analyses of the same cell, and does not permit arbitrary cell subpopulation sorting. Conventional microfluidic devices rely on cell-specific antibodies for cell selection, wherein the antibodies that are bound to the microfluidic device substrate selectively bind to cells expressing the desired antigen. Conventional microfluidic devices fail to allow for subsequent cell removal without cell damage, and only capture the cells expressing the specific antigen; non-expressing cells, which could also be desired, are not captured by these systems. Cellular filters can separate sample components based on size without significant cell damage, but suffer from clogging and do not allow for specific cell identification, isolation, and retrieval.\nThus, there is a need in the cell sorting field to create a new and useful cell capture and analysis system."}
{"text": "The subject invention relates to smoking articles, cigarettes for example, and smoking material therefor.\nThe patents literature contains many proposals for smoking materials for use in place of conventional cut tobacco cigarette filler.\nIt is an object of the subject invention to provide new smoking materials which provide for mainstream smoke which although containing low levels of tobacco derived components, is fully acceptable to the consumer.\nThe subject invention provides a smoking material comprising a non-polyol aerosol generator, up to 20% by weight tobacco, binder at not more than 20% by weight and not less than 30% by weight inorganic filler.\nTriethylene glycol diacetate (‘TEGDA’), glycerol triacetate (‘triacetin’) or glycerol diacetate (‘diacetin’), for example, can be used as the non-polyol aerosol generator either individually or in combination. As is well known to those skilled in smoking article science and technology, TEGDA and triacetin are substances with an established use as plasticisers (bonding agents) for cellulose acetate cigarette-filter tow. It was thus surprising to find that these substances and similar substances, when used as aerosol generating substances in smoking materials according to the subject invention, provide aerosols in mainstream smoke which smokers register as very acceptable.\nNon-polyol aerosol generator is usefully present in smoking materials of the subject invention at levels in a range of about 2% to about 30% by weight. More usually the range will be about 5% to about 20% by weight, and more usually 5-15%.\nTotal aerosol generator can include a polyol aerosol generator, such as for example one or more of glycerol, propylene glycol and triethylene glycol.\nSmoking materials according to the subject invention can, as will be readily appreciated by those skilled in the art, be fabricated by slurrying the components, in fine particulate form, with water and casting the slurry to sheet form on a band or wire sheet-forming machine or on a heated drum. An alternative is to feed a mixture of the components, together with water, to an extruder. The product of casting or extrusion is suitably cut and shredded to provide smoking material of particulate form.\nIf the components used to provide the smoking material do not include tobacco, then advantageously the smoking material, in particulate form, is blended with particulate tobacco. In such case, the smoking material expediently accounts for at least about 30% by weight of the blend. Suitably, the smoking material will account for the majority, by weight, of the blend, that is more than 50% by weight of the blend. The tobacco in the blend may be expanded tobacco.\nA class of substance suitable for the selection therefrom of binder in smoking materials according to the subject invention is the alginates. Sodium alginate has been found to be advantageous. Other suitable binder substances are celluloses or modified celluloses, hydroxypropyl cellulose or carboxymethyl cellulose, for example, starches or modified starches and natural gums.\nSuitable substances for use as inorganic filler are calcium carbonate, perlite, vermiculite, diatomaceous earth, colloidal silica, magnesium oxide, magnesium sulphate, magnesium carbonate or other low density inorganic filler materials known to those skilled in the art.\nIn smoking materials according to the subject invention inorganic filler is preferably present at a level of at least about 35%, preferably at least about 40% and more preferably at least about 45% by weight.\nIn smoking materials according to the subject invention tobacco may be present at less than 15%, preferably less than 10% and more preferably less than 5%. The binder of smoking materials according to the invention may be present at less than 15%, and more preferably less than 10% by weight.\nSmoking materials according to the subject invention may comprise one or more mechanical stabiliser or strengthening materials, examples being cocoa, sugar and fibre, paper fibre for instance. Expansion medium, such as starch, pullulan or other polysaccharides or foaming agents, for example, and high fat or high oil materials, such as cocoa butter or olive oil, corn oil, for example, may also be advantageously included.\nSmoking material according to the invention may be used in a conventional smoking article, either blended with another smoking material, which may be tobacco material, or not or in a coaxial arrangement, for example.\nAccording to a second aspect of the subject invention a smoking article comprising a smokable rod and a filter, said rod comprising a core and an outer part, said core comprising a first particulate smoking material and a first cigarette paper wrapper enwrapping said first smoking material and said outer part comprising a second particulate smoking material disposed annularly about said core and a second cigarette paper wrapper enwrapping said second smoking material, wherein either of said first or said second smoking material comprises a non-polyol aerosol generator, up to 20% by weight tobacco, binder at not more than 20% by weight and not less than 30% by weight inorganic filler.\nAs will be observed, one of the said smoking materials of the smoking article of the second aspect of the subject invention is as per the smoking material of the first above recited aspect of the subject invention. The smoking material of the first above recited aspect of the subject invention may be present in both the first smoking material and the second smoking material of the smoking article above.\nPreferably the said second smoking material of the smoking article is as per the smoking material of the first above recited aspect of the subject invention.\nExpediently, the first smoking material is cut tobacco filler, suitably cut lamina filler.\nMuch by preference, in smoking articles in accordance with the subject invention both the core and the outer part of the smokable rod extend over the full length of the rod. It is much by preference too for the core to be disposed coaxially of the rod.\nIn a third aspect of the invention, there is provided a smoking article comprising a smokable rod and a filter, the smokable rod comprising smoking material and the filter comprising filtration material provided with an elutable aerosol generator which is a polyol and/or a non-polyol aerosol generator, the aerosol generator being elutable from the filter upon smoking of the smoking article to enhance the aerosol provided from the smoking material.\nThe smoking article may be of a coaxial arrangement or non-coaxial arrangement.\nAdvantageously the smoking material comprises an aerosol generator, a binder and an inorganic filler, and is more advantageously comprised of a non-polyol aerosol generator, up to 20% by weight tobacco, binder at not more than 20% by weight and not less than 30% by weight inorganic filler. Alternatively the smoking material may comprise one or more of the aerosol generating means or aerosol generating fuel sources of our co-pending application filed on 6 Sep. 1995 under International Patent Application No. PCT/GB 95/02110. As a further alternative, the smoking material may comprise conventional cut tobacco leaf material.\nThe aerosol generator may be a polyhydric alcohol, an ester, a high boiling point hydrocarbon, glycerol, propylene glycol, triethylene glycol, methylene glycol, methyl citrate, triacetin or diacetin, either alone or in combination.\nIn a smoking article of an exterior circumference conventional for a cigarette, i.e. c.25 mm, the exterior circumference of the said core is suitably about 17 mm."}
{"text": "1. Field of the Invention\nThe present invention relates to a misting apparatus for a boat and more particularly pertains to creating a fine mist to be sprayed into a boat with a misting apparatus for a boat.\n2. Description of the Prior Art\nThe use of misting devices is known in the prior art. More specifically, misting devices heretofore devised and utilized for the purpose of delivering a water mist in numerous applications are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nBy way of example, U.S. Pat. No. 5,340,345 to Brodbeck et al. discloses a water pickup cooling apparatus for boat drive systems.\nU.S. Pat. No. 5,121,882 to Skidmore discloses a mist apparatus for sunbathing.\nU.S. Pat. No. 4,765,542 to Carlson discloses a liquid misting attachment for sunbather's chair.\nU.S. Pat. No. 3,759,449 to Ruthman et al. discloses a device for producing coolant mist.\nWhile these devices fulfill their respective, particular objective and requirements, the aforementioned patents do not describe a misting apparatus for a boat for creating a fine mist to be sprayed into a boat.\nIn this respect, the misting apparatus for a boat according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in doing so provides an apparatus primarily developed for the purpose of creating a fine mist to be sprayed into a boat.\nTherefore, it can be appreciated that there exists a continuing need for new and improved misting apparatus for a boat which can be used for creating a fine mist to be sprayed into a boat. In this regard, the present invention substantially fulfills this need."}
{"text": "The term, organic electronic device, refers to a device which requires electronic charge exchange between an electrode and an organic material by using holes and/or electrons. The organic electronic device can be largely classified into two types according to its operational principle as follows. One type is an electronic device having a configuration in which an exciton is formed in an organic material layer by photons flown from an external light source into the device and the exciton is separated into an electron and a hole, the electron and the hole formed are transported to a different electrode, respectively and used as a current source (voltage source), and the other type is an electric device having a configuration in which holes and/or electrons are injected into an organic material semiconductor forming an interface with an electrode by applying a voltage or current to two or more electrodes to allow the device to operate by means of the injected electron and hole.\nExamples of the organic electronic device comprise an organic light emitting device, an organic solar cell, an organic photoconductor (OPC) and an organic transistor, which all require a hole injecting or transporting material, an electron injecting or transporting material, or a light emitting material for driving the device.\nHereinafter, the organic light emitting device will be mainly and specifically described, but in the above-mentioned organic electronic devices, the hole injecting or transporting material, the electron injecting or transporting material, or the light emitting material functions according to a similar principle.\nIn general, the term “organic light emitting phenomenon” refers to a phenomenon in which electric energy is converted to light energy by means of an organic material. The organic light emitting device using the organic light emitting phenomenon has a structure usually comprising an anode, a cathode and an organic material layer interposed therebetween. Herein, the organic material layer may be mostly formed in a multilayer structure comprising layers of different materials, for example, the hole injecting layer, the hole transporting layer, the light emitting layer, the electron transporting layer, the electron injecting layer and the like, in order to improve efficiency and stability of the organic light emitting device. In the organic light emitting device having such a structure, when a voltage is applied between two electrodes, holes from the anode and electrons from a cathode are injected into the organic material layer, the holes and the electrons injected are combined together to form excitons. Further, when the excitons drop to a ground state, lights are emitted. Such the organic light emitting device is known to have characteristics such as self-luminescence, high brightness, high efficiency, low drive voltage, wide viewing angle, high contrast and high-speed response.\nThe materials used for the organic material layer of the organic light emitting device can be classified into a light emitting material and a charge-transporting material, for example, a hole injecting material, a hole transporting material, an electron transporting material and an electron injecting material, according to their functions. The light emitting materials can be divided into a high molecule light emitting material and a low molecule light emitting material according to the molecular weight thereof, and can be divided into a fluorescent material that is derived from a singlet excite state of electrons and a phosphorescence material that is derived from a triplet excite state of electrons according to the light emitting mechanism. In addition, the light emitting materials can be divided into a blue, green or red light emitting material and a yellow or orange light emitting material required for giving more natural color, according to a light emitting color.\nMeanwhile, in the case of when only one material is used as the light emitting material, the maximum light emitting wavelength is moved into the long wavelength by the interaction between molecules, the color purity is reduced, or efficiency of the device is lowered by the light emitting reducing effect. Accordingly, a host/dopant system can be used as the light emitting material for the purpose of enhancing the color purity and the light emitting efficiency through energy transfer. It is based on the principle that if a small amount of a dopant having a smaller energy gap than a host which forms a light emitting layer is mixed with the light emitting layer, excitons which are generated in the light emitting layer are transported to the dopant, thus emitting a light having a high efficiency. Here, since the wavelength of the host is moved according to the wavelength of the dopant, a light having a desired wavelength can be obtained according the kind of the dopant.\nIn order to allow the organic light emitting device to fully exhibit the above-mentioned excellent characteristics, a material constituting the organic material layer in the device, for example, a hole injecting material, a hole transporting material, a light emitting material, an electron transporting material and an electron injecting material should be essentially composed of a stable and efficient material. However, the development of a stable and efficient organic material layer material for the organic light emitting device has not yet been fully realized. Accordingly, the development of new materials is continuously desired. The need to develop the above material is the same as the case of the above-mentioned other organic electronic devices."}
{"text": "Spun yarns of synthetic staple fibers have been produced by cutting continuous filaments into staple fibers, which are then assembled into yarn in the same manner as fibers of cotton or wool. A simpler direct spinning process is also used wherein parallel continuous filaments are stretch-broken and drafted between input rolls and delivery rolls in a drafting zone to form a sliver of discontinuous fibers which is thereafter twisted to form a spun yarn as disclosed, for example, in New U.S. Pat. No. 2,721,440 or Preston U.S. Pat. No. 2,784,458. A true twisting operation is inherently slow. Higher speeds can be obtained by using a false-twisting device, as disclosed in Tissot et al. U.S. Pat. No. 2,946,181 or Field, Jr. U.S. Pat. No. 3,079,746. Tissot et al. teaches that the use of gelatin size will provide adequate strength for weaving or knitting. However, the yarn lacks strength after the size is removed during fabric finishing. Field, Jr., teaches the use of a twisting jet to produce a yarn having surface fibers wrapped about a core bundle which is substantially free of true twist. However, improvements in this product are desirable to make it even more like a conventional spun yarn.\nBunting, Jr., et al. U.S. Pat. No. 3,110,151 discloses at column 23, lines 33-52, a process in which staple fibers averaging 7.6 centimeters long and 2 denier per filament were fed as a roving to drafting rolls and then through an entangling jet device, referred to therein as an interlacing jet, for entangling into yarn. Such products are not comparable to conventional spun yarns in strength, cleanness (freedom from neps and slubs), and uniformity."}
{"text": "1. Field of the Invention\nThe present invention relates to lawn ornaments and more particularly pertains to a new decorative lawn ornament for providing inspirational three-dimensional ornamentation.\n2. Description of the Prior Art\nThe use of lawn ornaments is known in the prior art. More specifically, lawn ornaments heretofore devised and utilized are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nKnown prior art includes U.S. Pat. No. 2,646,497; U.S. Pat. No. Des. 373,097; U.S. Pat. No. Des. 371,987; U.S. Pat. No. Des. 304,347; U.S. Pat. No. 4,601,924; and U.S. Pat. No. 5,584,571.\nWhile these devices fulfill their respective, particular objectives and requirements, the aforementioned patents do not disclose a new decorative lawn ornament. The inventive device includes a main body frame formed in the shape of an angel and having a plurality of elongate frame members, a plurality of main body lights, a secondary frame formed in the shape of a bird, a plurality of secondary lights coupled to the secondary frame, a wing motor interior to the secondary frame, a rod motor interior to the main body frame, a rod, and a linkage assembly.\nIn these respects, the decorative lawn ornament according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in so doing provides an apparatus primarily developed for the purpose of providing inspirational three-dimensional ornamentation"}
{"text": "Semiconductor device manufacturing processes include various processes such as a lithographic process, an etching process, and an ion implantation process. After the end of each process and before shifting to the next process, a cleaning process and a drying process are performed so as to remove impurities and residues remaining on the surface of a semiconductor substrate to clean the surface of the semiconductor substrate.\nIn recent years, following the downscaling of elements, the aspect ratio of patterns on a semiconductor substrate has become higher. At a higher aspect ratio, there occurs a problem that capillary causes collapsing of the patterns on the semiconductor substrate in the drying process.\nTo deal with such a problem, there has been proposed a technique for making surfaces of the patterns on the semiconductor substrate water repellent and lowering the capillary that acts between the patterns and a chemical liquid or the DIW (Deionized Water). However, a water-repellent agent used for making the surface of the semiconductor substrate water repellent is often deactivated after reacting to the water. For example, it often occurs in a cleaning device that the water-repellent agent is deactivated after reacting to the water in a chamber. If such deactivation of the water-repellent agent occurs, the water-repellent agent is unable to make the surface of the semiconductor substrate water repellent and to suppress collapsing of the patterns on the semiconductor substrate resulting from the capillary. Furthermore, if the water-repellent agent is deactivated while spreading throughout the surface of the semiconductor substrate, water repellency varies on the surface of the semiconductor substrate."}
{"text": "PTFE (polytetrafluoroethylene) is used as a material for packings for sealing an annular gap between a shaft and a housing that move relative to each other. PTFE can be used under high temperatures and high pressures, and has excellent characteristics such as high chemical resistance and excellent sliding properties. On the other hand, PTFE has a shortcoming that it is prone to plastic deformation (settling) particularly when used in high temperature environments. If settling occurs in a sealing lip of the packing, it leads to deterioration of sealing properties.\nA technique for maintaining the sealing lip in tight contact with an outer circumferential surface of a shaft or an inner circumferential surface of a shaft hole of a housing is known (see Patent Literatures 1 and 2), wherein a metal spring member is mounted in an annular mounting groove formed in the packing.\nHowever, PTFE has a higher linear expansion coefficient than that of metal materials commonly used as the housing material. In high temperature environments, therefore, the sealing lip on an outer circumferential side of the packing (outer circumferential lip) tends to deform so as to expand radially outward, so that it is firmly pressed against the inner circumferential surface of the shaft hole of the housing. As a result, a large settling occurs in the outer circumferential lip, hence, as a countermeasure for settling, it is not sufficient to simply mount a spring member.\nConventionally, in order to prevent the metal spring member from coming off, an annular protrusion that protrudes radially inward is provided at the distal end of the outer circumferential lip. The outer circumferential lip thus has a relatively high rigidity."}
{"text": "The invention relates to a process for the preparation of an olefin polymerization catalyst component containing magnesium, titanium, halogen and an electron donor. The invention also relates to such a catalyst component and its use for the polymerization of xcex1-olefins such as propene.\nGenerally, so called Ziegler-Natta catalyst components of the above kind have been prepared by reacting a magnesium halide-alcohol complex with a titanium tetra-halide and an electron donor which usually is a phthalic acid diester. The preparation involves the use of large amounts of reagents and washing liquids, which are difficult to handle. Additionally, byproducts are formed, which cannot easily be regenerated or destroyed, but form an environmental problem.\nFor example, the preparation of a conventional polypropene catalyst component involves the reaction of a magnesium dichloride-alcohol complex with titanium tetrachloride to give reactive xcex2-magnesium dichloride as intermediate and hydrogen chloride and titanium alkoxy trichloride as byproducts. Then, the reactive xcex2-magnesium dichloride intermediate is activated with further titanium tetrachloride to give said catalyst component (the treatment with a titanium halide such as titanium tetrachloride is henceforth called titanation).\nThe titanium alkoxy trichloride byproduct formed in the titanation is a catalyst poison and must be carefully removed by extensive washing using large amounts of titanium tetrachloride. Further, the titanium alkoxy trichloride must be carefully separated from the titanium tetrachloride washing liquid, if the latter is to be reused e.g. for activating the reactive xcex2-magnesium dichloride. Finally, the titanium alkoxy trichloride is a problem waste, which is difficult to dispose of.\nThus, in a typical propene polymerization catalyst component preparation involving two titanations and three heptane washes, one mol of produced catalyst component (mol Mg) requires about 40 mol of titanium tetrachloride e.g. as washing liquid to be circulated (see Table 15 below), and produces as problem waste about three mol of titanium alkoxy trichloride as well as about three mol of hydrogen chloride.\nSumitomo, EP 0 748 820 A1 (hereinafter referred to as xe2x80x9cSumitomoxe2x80x9d), has prepared dialkoxy magnesium, reacted it with titanium tetrachloride to form an intermediate and then reacted the intermediate with phthalic acid dichloride to form a catalytically active propene polymerization catalyst component. The activity was raised by repeated titanations, as well as repeated washes with toluene and hexane. See page 10, lines 14 to 37, of said publication.\nSaid process of Sumitomo has avoided the reaction between the magnesium dichloride-alcohol complex and titanium tetrachloride, and thereby eliminated the formation of catalytically poisonous titanium alkoxy trichloride byproduct. However, as much as four titanations and hydrocarbon treatments are still needed to give satisfactory catalytic activity.\nThe purpose of the present invention is to provide a process which results in a catalyst component having satisfactory activity without producing harmful byproducts such as said titanium alkoxy trichloride or requiring the use of a large amounts of titanation reagent and/or washing liquid.\nThe problem described above has now been solved with a novel process for the preparation of a catalyst component of the above type, which is mainly characterized by the steps of:\n(i) reacting a titaniumless magnesium compound (a) containing an alkoxy moiety, which titaniumless magnesium compound is selected from the group consisting of a compound or conplex containing halogen and alkoxide linked to magnesium, a complex containing a magnesium dihalide and an alcohol, and a non-complex magnesium dialkoxide, with a halogen compound (b) being capable of forming the electron donor by replacement of its halogen by said alkoxy moiety, to give an intermediate (ab); and\n(ii) reacting said intermediate (ab) with a titanium halide (c), or\n(i)xe2x80x2 reacting a titaniumless magnesium compound (a) containing an alkoxy moiety, which titaniumless magnesium compound is selected from the group consisting of a compound or complex containing halogen and alkoxide linked to magnesium, and a complex containing a magnesium dihalide and an alcohol, with a titanium halide (c), to give an intermediate (ac), and\n(ii)xe2x80x2 reacting said intermediate (ac) with a halogen compound (b) being capable of forming the electron donor by replacement of its halogen by said alkoxy moiety.\nIt was found by the applicant, that the activity of a stoichiometric catalyst component, comprising a magnesium dihalide, a titanium tetrahalide and an electron donor, is the higher, the more magnesium dihalide it contains. Thus it is believed, that the purpose of the repeated toluene washes of e.g. Sumitomo has partly been to remove titanium tetrachloride and electron donor from the catalyst component precursor in order to raise the magnesium dichloride content and thus the catalytic activity, of the final catalyst component. The present invention solves the problem in another way. In the claimed process, magnesium dihalide is included or synthesized as part of the reacting material before any titanation takes place, and thus, the need for repeated cycles of titanation and washing is significantly reduced.\nPreferably one, most preferably all of steps (i), (ii), (i)xe2x80x2 and (ii)xe2x80x2 are performed in solution. Then, the reaction product of step (ii) or (ii)xe2x80x2 is preferably recovered by precipitation.\nAccording to one embodiment of the present invention, said compounds (a), (b) and (c) are in the claimed process contacted in essentially stoichiometric amounts. According to another embodiment, a stoichiometric excess, preferably a 5-20 fold excess, of said titanium halide (c) with respect to said magnesium compound (a), gives even better results.\nSaid halogen compound (b) used in the claimed process is an electron donor precursor, i.e. itself capable of forming the electron donor of the catalyst component by replacement of its halogen by an alkoxy group. By electron donor is in this connection meant an electron donor which forms a part of the titanous catalyst component produced by the claimed process and is in the art also called an internal electron donor. Such halogen compounds (b) are, e.g., C1-C20 alkyl halides, C7-C27 aralkyl halides and C2-C22 acyl halides, which react with alkoxy compounds to replace their halogen with the alkoxy group of the alkoxy compound and form e.g. the corresponding ethers and esters acting as internal electron donors.\nPreferably, said halogen compound (b) is an organic acid halide having the formula Rxe2x80x3(COXxe2x80x2)n, wherein Rxe2x80x3 is an n-valent organic group having 1-20 carbon atoms, preferably an n-valent benzene ring, Xxe2x80x2 is a halogen, preferably chlorine; and n is the valence of Rxe2x80x3 and is an integer 1 to 6, preferably 1, 2, 3 or 4, more preferably 2. Most preferably, said halogen compound is phthalic acid dichloride Ph(COCl)2, wherein Ph is o-phenylene. Correspondingly, the electron donor formed therefrom is preferably an organic acid ester having the formula Rxe2x80x3(COOR)n, wherein R is an n-valent C1-C20 aliphatic group or an n-valent C7-C27 araliphatic group and Rxe2x80x3 and n are the same as above, and more preferably a phthalic acid diester Ph(COOR)2, wherein R is a C1-C20 alkyl or a C7-C27 aralkyl, more preferably a C1-C16 alkyl. Most preferably said electron donor is dioctyl phthalate.\nThe titanium halide (c) used in the claimed process is preferably a titanium halide of the formula (ORxe2x80x2)kTiX4-k, wherein Rxe2x80x2 is an alkyl group having 1 to 10 carbon atoms or an aralkyl group having 7 to 16 carbon atoms, X is a halogen and k is 0 to 3. More preferably, said titanium halide (c) is a titanium tetrahalide TiX4, wherein X is the same as above, most preferably titanium tetrachloride TiCl4.\nIt is preferable, if said titaniumless magnesium compound (a) is not a part of a solid magnesium halide, e.g. in the form of complex molecules on the surface of a solid magnesium halide carrier, but form a separate compound with an essentially stoichiometric composition. Often, said titaniumless magnesium compound is a complex. A complex is, according to Rxc3x6mpps Chernie-Lexicon, 7. Edition, Franckh\"\"sche Verlagshandlung, W. Keller and Co., Stuttgart, 1973, page 1831, xe2x80x9ca derived name of compounds of higher order, which originate from the combination of molecules,xe2x80x94unlike compounds of first order, in the creation of which atoms participatexe2x80x9d.\nAccording to one embodiment of the invention, the titaniumless magnesium compound (a) used in the claimed process is a titaniumless complex of the formula [MgX2]x.[K(OR)m]y, wherein X is a halogen, K is hydrogen, a metal of group 1, 2 or 13 of the Periodic Table, R is an alkyl having 1 to 20 carbon atoms, an aralkyl having 7 to 27 carbon atoms or an acyl having 2 to 22 carbon atoms, x is 0 to 20, m is the valence of K and is an integer from 1 to 6, and y is 1 to 20.\nThe magnesium dihalide MgX2 of said titaniumless complex (a) can be selected from magnesium chloride, magnesium bromide and magnesium iodide. Preferably, it is magnesium dichloride.\nThe alkoxy compound K(OR)m is in its most general form defined as a component of said titaniumless complex (a), which complex (a) is reacted further with said halogen compound (b) and said titanium halide (c) or said titanium halide (c) and said halogen compound (b). The alkoxy compound is, however, more closely defined in the following description of three main embodiments of the claimed process.\nThe gist of the invention is to choose the reactants and their order of reaction so that the magnesium dihalide is present when the titanium halide (c) is reacted."}
{"text": "1. Technical Field\nThe present invention relates to a charging unit. In particular, although not exclusively, the present invention relates to a charging unit comprising a switched mode power supply.\n2. Description of Related Art\nPortable devices are often powered by a rechargeable battery. Most commonly, the rechargeable battery is charged by a charger including an AC to DC converter and being connected to an AC mains power supply. In addition, DC chargers such as vehicle chargers can be used to recharge batteries. Such chargers typically can provide charging currents of around 1 A.\nSome portable devices have a USB port for connection to a host device. The USB connection can be used for data transfer. For example, a mobile phone may be connected to a computer for synchronising their respective address books.\nThe USB specification also provides that a USB connection is able to supply power to a device. However, the power that can be supplied is limited. The USB specification recites a DC voltage of 4.75 to 5.25V, and the current that can be provided through a USB connection is limited to around 0.5 A. Thus, the time taken to charge a portable device using only USB charging can be significantly longer than is found when using AC mains charging.\nPortable devices that allow both USB charging and AC mains charging have been proposed. One example is described in US 2005/0141208. This document describes a USB device that can be connected simultaneously to an AC adapter and to a USB port on a host device by a USB connection. When the AC adapter is connected, the USB device operates as a self-powered USB device, i.e. it is not supplied with power from the host device."}
{"text": "Digital Subscriber Line (DSL) has become popular in recent years because it can provide high-bandwidth connections between a telephone switching station and a home or office over existing telephone lines. DSL technologies use sophisticated modulation schemes to pack data onto copper wires of plain old telephone system (POTS). Among the exiting DSL technologies, Asymmetric Digital Subscriber Line (ADSL) is the most practical since it fits average users' need for greater downstream bandwidth than upstream bandwidth.\nDiscrete Multi-tone (DMT) is a method of modulating a DSL signal for transmission on a usable frequency spectrum divided into a plurality of sub-carriers or frequency channels. For example, a spectrum may be divided into 256 channels (“bins”) of 4.3125 kHz each. The center frequency of a bin is 4.3125 kHz multiplied by the bin number. Thus, a specific bin may be referred to by its bin number instead of its actual frequency range. Dividing the frequency spectrum into multiple channels reduces cross-talk in signal transmissions. In a DMT system, a usable frequency spectrum is typically allocated to the upstream and downstream transmissions based on a standard scheme. For example, according to International Telecommunication Union (ITU) standard G.992.1 (Annex A), the downstream transmission may occupy bins 6 through 255 and the upstream transmission may occupy bins 6 through 31. In a common operation mode, to reduce near-end cross-talk (NEXT), the downstream transmission may start from bin 33. Other standards or schemes also exist that allow an overlap in bin allocations to the upstream and downstream traffic.\nHowever, none of the fixed bin allocation schemes can consistently achieve optimal data rates. Due to the ever-changing operation environment (e.g., noise level, loop attenuation and echo rejection rate), quality of a connection based on a particular scheme usually drifts over time. Especially when the cross-talk noises become highly asymmetric, it may be difficult to rely on a fixed bin allocation scheme to maintain link performance. As a result, the available bandwidth is not fully utilized and customers may become dissatisfied with inconsistent link performance.\nIn view of the foregoing, it would be desirable to provide an efficient and cost effective solution for bin allocation that overcomes the above-described inadequacies and shortcomings."}
{"text": "The prior art is replete with light duty liquid detergent compositions containing nonionic surfactants in combination with anionic and/or betaine surfactants wherein the nonionic detergent is not the major active surfactant. In U.S. Pat. No. 3,658,985 an anionic based shampoo contains a minor amount of a fatty acid alkanolamide. U.S. Pat. No. 3,769,398 discloses a betaine-based shampoo containing minor amounts of nonionic surfactants. This patent states that the low foaming properties of nonionic detergents renders its use in shampoo compositions non-preferred. U.S. Pat. No. 4,329,335 also discloses a shampoo containing a betaine surfactant as the major ingredient and minor amounts of a nonionic surfactant and of a fatty acid mono- or di-ethanolamide. U.S. Pat. No. 4,259,204 discloses a shampoo comprising 0.8 to 20% by weight of an anionic phosphoric acid ester and one additional surfactant which may be either anionic, amphoteric, or nonionic. U.S. Pat. No. 4,329,334 discloses an anionic-amphoteric based shampoo containing a major amount of anionic surfactant and lesser amounts of a betaine and nonionic surfactants.\nU.S. Pat. No. 3,935,129 discloses a liquid cleaning composition containing an alkali metal silicate, urea, glycerin, triethanolamine, an anionic detergent and a nonionic detergent. The silicate content determines the amount of anionic and/or nonionic detergent in the liquid cleaning composition. However, the foaming properties of these detergent compositions are not discussed therein.\nU.S. Pat. No. 4,129,515 discloses a heavy duty liquid detergent for laundering fabrics comprising a mixture of substantially equal amounts of anionic and nonionic surfactants, alkanolamines and magnesium salts, and, optionally, zwitterionic surfactants as suds modifiers.\nU.S. Pat. No. 4,224,195 discloses an aqueous detergent composition for laundering socks or stockings comprising a specific group of nonionic detergents, namely, an ethylene oxide of a secondary alcohol, a specific group of anionic detergents, namely, a sulfuric ester salt of an ethylene oxide adduct of a secondary alcohol, and an amphoteric surfactant which may be a betaine, wherein either the anionic or nonionic surfactant may be the major ingredient.\nThe prior art also discloses detergent compositions containing all nonionic surfactants as shown in U.S. Pat. Nos. 4,154,706 and 4,329,336 wherein the shampoo compositions contain a plurality of particular nonionic surfactants in order to affect desirable foaming and detersive properties despite the fact that nonionic surfactants are usually deficient in such properties.\nU.S. Pat. No. 4,013,787 discloses a piperazine based polymer in conditioning and shampoo compositions which may contain all nonionic surfactant or all anionic surfactant.\nU.S. Pat. No. 4,450,091 discloses high viscosity shampoo compositions containing a blend of an amphoteric betaine surfactant, a polyoxybutylenepolyoxyethylene nonionic detergent, an anionic surfactant, a fatty acid alkanolamide and a polyoxyalkylene glycol fatty ester. But, none of the exemplified compositions contain an active ingredient mixture wherein the nonionic detergent is present in major proportion which is probably due to the low foaming properties of the polyoxybutylene polyoxyethylene nonionic detergent.\nU.S. Pat. No. 4,595,526 describes a composition comprising a nonionic surfactant, a betaine surfactant, an anionic surfactant and a C.sub.12 -C.sub.14 fatty acid monoethanolamide foam stabilizer."}
{"text": "Conventionally, a signal processing system has a TI-ADC (time-interleaved analogue digital converter), an ADC (analogue-to-digital converter), and a correction processor. The correction processor estimates an error to be corrected by using the steepest descent method of the covariance norm and corrects an error caused by mismatch between interleaves by using the estimation result and an adaptive filter.\nHowever, since the adaptive filter is large-scale, the size of the correction processor is extremely large, and a process amount necessary to correct an error caused by mismatch between interleaves is extremely large. As a result, it is difficult to realize small size and high speed of the TI-ADC."}
{"text": "1. Field of the Invention\nThe present invention relates to integrated circuit fabrication, and more particularly to methods for fabricating high-density integrated circuit devices.\n2. Description of Related Art\nPhotolithographic processes can be used to form a variety of integrated circuit structures on a semiconductor wafer. In photolithography, features of these structures are typically created by exposing a mask pattern (or reticle) to project an image onto a wafer that is coated with light sensitive material such as photo resist. After exposure, the pattern formed in the photo resist may then be transferred to an underlying layer (e.g. metal, polysilicon, etc.) through etching, thereby creating the desired features.\nOne problem associated with manufacturing devices having very small features arises because of variations introduced by the photolithographic processes. Specifically, resist material properties, process conditions, optical distortions and other factors can cause systematic and random deviations in the etched shapes of the features from their desired shapes. Examples of deviations include corner-rounding, line-shortening and line edge roughness.\nIn a typical lithographic patterning process, a line of resist is used as an etch mask to create a corresponding line of material in the underlying layer. In such a case, the deviations in the patterned line of resist will be transferred to the critical dimensions of the etched line in the underlying layer. As process technologies continue to shrink, these deviations become a greater percentage of the critical dimension of the etched lines, which can reduce yield and result in significant performance variability in devices such as transistors implemented utilizing these etched lines.\nAccordingly, it is desirable to provide high-density structures such as integrated circuit devices which overcome or alleviate issues caused by deviations introduced by photolithographic processes, thereby improving performance and manufacturing yield of such devices."}
{"text": "The present invention relates generally to broadband transmission of information and, more particularly, to time synchronization of a received OFDM signal.\nOrthogonal Frequency Division Multiplexing (OFDM) is a spread spectrum technology wherein the available bandwidth is subdivided into a number of discrete channels or subcarriers that are overlapping and orthogonal to each other. Each channel has a well defined frequency. Data are transmitted in the form of symbols that encompass some number of subcarrier frequencies. The amplitude and/or phase of the subcarriers represents the encoded information. Before a received OFDM symbol can be decoded, frequency and time synchronization of the signal have to be determined. This process is referred to as synchronization.\nTime synchronization may be achieved through the use of time-based correlation in the time domain or phase correlation in the frequency domain. Phase correlation has several advantages over time-based correlation. Phase correlation is more robust as it is significantly less affected by amplitude variations than time-based correlation. Also, symbols can be synchronized in the presence of a much stronger jammer by correlation in the frequency domain.\nOFDM symbols can be synchronized in the frequency domain by correlating the phases of the carriers with the phases of a reference symbol. Alternatively, or in addition, autocorrelation can be performed by correlating the phases of a received symbol with those of an earlier received symbol.\nTime synchronization of the received OFDM signals is more difficult to achieve, in particular, in the presence of narrow band interference in the OFDM channel. Apodizing window functions, such as Hanning windows, are typically used to deal with the effects of a narrowband interferer, such as a jammer signal, in samples that are passed to a phase correlator. While the apodizing window confines interference from the narrowband interferer to only several of the OFDM subcarriers, it introduces intercarrier interference into the subcarrier frequencies. The intercarrier interference can result in destructive canceling (i.e., frequency nulls\") in transmitting subcarriers of equal amplitude and phase. In the presence of a frequency null, a phase correlator may miss a peak correlation pulse and thus provide unreliable synchronization. The intercarrier interference can also cause phase and amplitude distortion in each of the transmitted subcarriers processed by the apodizing window. Phase distortion in sample phases, as well as in the autocorrelation or reference phases with which they are compared, can result in poor phase correlation. Even using a best possible set of reference phases (typically, nonwindowed referenced phases), in conjunction with Hanning windowed phases reduces correlation by some factor. That is, it takes a much stronger signal to perform correlation than would be needed if a Hanning window was not used."}
{"text": "This invention relates to a method of and apparatus for determining the no-flow temperature of a liquid, such as domestic heating oil, plastic, paint, ice cream, margarine, oil, distillate fuel and aviation fuel. The determination of the no-flow temperature of such materials is important, because their production and handling is limited by their low temperature properties. It is evident that significant problems may arise, for example, if fuels such as aviation fuel or domestic fuel fail to flow through the handling system to a burner because of their limiting low temperature properties.\nOne of the critical specifications of liquids such as jet fuel, is the freeze point which is defined as the temperature at which wax crystals formed on cooling disappear when the temperature of an agitated fuel sample is allowed to rise. A major shortcoming of the freeze point specification is that it bears little relation to the fuel property that limits low temperature performance, viz., the ability of the fuel to flow from a fuel tank at low temperature. A prior art cold flow test consists of characterizing fuel by its hold-up factor, which is defined as the amount of fuel that is unable to flow from a first section of a tester to a second section and varies with the temperature of the fuel. The most important test parameter is the \"zero percent hold-up temperature\", which represents the lowest temperature to which the fuel can be subjected before total hold-up will occur, that is, the point at which none of the fuel can flow. It is this temperature that is likely to define the limit of a specification based on the low temperature flow properties of the fuel. The zero percent hold-up temperature cannot be determined directly by the prior art techniques. Rather, the usual procedure consists of constructing a graph of the temperature dependence of the hold-up factor over a range of readily determined values, which are typically between 10-90% hold-up, and extrapolating the plot to zero percent hold-up to obtain the zero percent hold-up temperature by this method, several points are necessary to construct the hold-up temperature profile, particularly in the region of low hold-up factors. Obviously, this is a time-consuming exercise.\nTherefore, it is an object of the present invention to provide a method of and an apparatus for determining the no-flow temperature of liquids."}
{"text": "A number of ink-jet ink sets utilize water-soluble dyes, which are highly soluble in water, in order to achieve the right color gamut. Examples of patents that are directed to ink-jet ink sets for achieving the desired color gamut include U.S. Pat. No. 5,108,504, issued Apr. 28, 1992, to Loren E. Johnson et al and entitled \"High Chroma Color Dye-Set for Use in Ink-Jet Inks Employing Polysaccharide-Containing Vehicles\"; U.S. Pat. No. 5,143,547, issued Sep. 1, 1992, to William D. Kappele and entitled \"Specific Dye Set for Thermal Ink-Jet Printing on Plain and Coated Papers\"; U.S. Pat. No. 5,145,519, issued Sep. 8, 1992, to William D. Kappele and entitled \"Specific Dye Set for Thermal Ink-Jet Printing\"; U.S. Pat. No. 5,185,034, issued Feb. 9, 1993 and entitled \"Ink-Jet Inks with Improved Colors and Plain Paper Capacity\"; and U.S. Pat. No. 5,534,051, issued Jul. 9, 1996, to Hiang P. Lauw and entitled \"Specific Dye Set for Thermal Ink-Jet Printing\"; all of the foregoing cited patents are assigned to the same assignee as the present application.\nHowever, many dye-based inks, after being printed, exhibit poor water-fastness. There is an increased demand by consumers on the permanence or durability of ink-jet print, including water-fastness, light-fastness, smear-fastness, smudge-fastness, etc. The purpose of the present invention is to improve the durability of water-soluble dye-based inks, especially the water-fastness, smear-fastness, smudge-fastness, and light-fastness, lumped together under the term \"color permanence\", of such inks."}
{"text": "Slidably mounted building products, such as sliding windows, single hung and double hung windows and sliding doors, are well known in the new construction and remodeling industry. These products are typically operated by a user exerting a sliding-type force on the movable sash in order to open or close the same. As will be appreciated, the amount of force that needs to be exerted can at times be substantial, and thus physically handicapped persons or those persons with inadequate strength may not be able to operate them.\nFurthermore, with the advent of so-called \"smart homes\" it is desired to be able to control electronically or by computer the opening and closing of windows and doors, without the need for a person to be present to open and close the same. Currently existing windows and doors with sliding sashes would be extremely difficult to convert to automatic operation.\nIn the first half of the 20th century, several patents were issued which disclosed crank type window operating apparatus for slidably mounted windows. For example, Lovell U.S. Pat. No. 1,152,425 shows a double hung window which can be moved up and down by use of a hand crank. The hand crank turns a bevel gear connected to a left-hand rod and a right-hand rod, the rods turning another bevel gear which in turn rotates a longitudinally disposed screw mounted in each of the vertical window frame members. A movable nut is threadably mounted on each screw and is also secured to the window sash. Rotation of the screws causes vertical movement of the movable nuts, which in turn lowers or lifts the window sash in the window frame. See also U.S. Pat. Nos. 1,571,819; 1,966,815; 2,545,449; 2,913,920; and 3,022,065.\nAlthough this prior art has been known for many years, there are today virtually no commercially available slidably mounted doors and windows which are operated by hand crank or motor driven mechanisms. There is, therefore, a need for a modern operating mechanism for slidably mounted building products that is inexpensive as well as easy to use. The mechanism must be useful for windows and doors having wood, metal or polymeric frames and must be reliable and easy to install in both new and existing homes."}
{"text": "With development of display technology, a display device with large size, high quality and low cost is the trend. Quality of a color filter which is a major component of a display device (for example, a TFT-LCD display device) determines the display effect of the display device.\nAt present, a proximity exposure machine optical system is mostly used to perform exposure when fabricating a color filter. Specifically, as shown in FIG. 1, the above mentioned optical system includes an exposure lamp 101, a first planar mirror 102, a fly eye lens 103, a concave mirror 104 and a second planar mirror 105, and further includes a mask plate 106, a substrate 107, an exposure region 108 and a photoresist (not shown in FIG. 1). In this optical system, the exposure lamp 101 is mainly used to emit a light beam to the first planar mirror 102, the first planar mirror 102 is mainly used to reflect an incident light beam onto the fly eye lens 103, the fly eye lens 103 is mainly used to split the received light beam into narrow light beams and then refract the narrow light beams onto the concave mirror 104, the concave mirror 104 is mainly used to reflect the received light beam onto the second planar mirror 105, and the second planar mirror 105 is mainly used to reflect the received light beam onto the mask plate 106. Generally, the concave mirror 104 and the second planar mirror 105 are designed based on a principle that reflectivities of light at respective positions on the mirror are equal so as to guarantee the uniformity of light at different positions in the optical system.\nHere, the fly eye lens 103 may include a first lens assembly and a second lens assembly. The first lens assembly includes a plurality of lenses forming a first lens face, and the second lens assembly includes a plurality of lenses forming a second lens face. The first lens face is used to split a broad light beam reflected from the first planar mirror 102 into narrow light beams and then refract the narrow light beams onto the second lens face, and the second lens face is used to dispersively refract the received narrow light beams onto the concave mirror 104. In general, there is positional correspondence existing between the fly eye lens 103 and the concave mirror 104, that is, a lens closer to the center on the second lens face contributes more illuminance to a point closer to the center on a surface of the concave mirror 104 but contributes less illuminance to a point farther from the center on a surface of the concave mirror 104; on the contrary, a lens farther from the center on the second lens face contributes more illuminance to a point farther from the center on the surface of the concave mirror 104 but contributes less illuminance to a point closer to the center on the surface of the concave mirror 104.\nWhen performing the proximity exposure, generally, the illuminance of incident light on the mask plate is evenly distributed, that is, illuminance at any region on the mask plate is identical. In this case, as shown in FIG. 2, when fabricating a color filter using a substrate 201 (for example, a glass substrate), a mask plate 202 and a photoresist 203, a distance between the mask plate 202 before being bent and the substrate 201 ranges from one hundred micrometers to several hundred micrometers. As such, as to the mask plate, it is bent when performing the exposure, which causes the exposure distance in a vertical direction between the substrate 201 and the central region of the mask plate 202 to differ from that between the substrate 201 and an edge region of the mask plate 202. For example, as shown in FIG. 2, the exposure distance G2 between the substrate 201 and the central region of the mask plate 202 is smaller than the exposure distance G1 between the substrate 201 and the edge region of the mask plate 202. During a usual fabricating process, since the proximity exposure machine optical system is influenced by factors such as diffraction angle and light parallelism when performing the exposure, the critical dimension (CD) of the pattern exposed on the substrate 201 is increased as the exposure distance of each region is increased (as shown in FIG. 2). Thus, in a case where the illuminance is fixed, since the exposure distance in the central region is shorter than that in an edge region, the critical dimension of the pattern exposed on the central region is smaller than that of the pattern exposed on an edge region, that is, the critical dimensions of the exposed patterns are not uniform (for example, CD2>CD1, as shown in FIG. 1). In other words, the critical dimensions of the patterns exposed by the above-mentioned optical system have a poor uniformity.\nIt can be seen from the above fabricating manner that when the proximity exposure machine optical system performs the exposure, the bending of the mask plate causes exposure distances in respective regions to vary so that the critical dimensions of the patterns exposed on the substrate have a poor uniformity and the quality of the color filter is further affected."}
{"text": "1. Field of the Invention\nThe invention relates to a charging apparatus for an electric storage device, a vehicle equipped with the charging apparatus, and a method of controlling the charging apparatus.\n2. Description of Related Art\nIn recent years, vehicles that run with driving force using electric power stored in an electric storage device installed on the vehicle have been attracting attention, as environment-conscious vehicles. Examples of the electric storage device include a second battery and a capacitor. The vehicles of the above type include, for example, electric vehicles, hybrid vehicles, fuel-cell vehicles, etc. It has been proposed to charge the electric storage device installed on the above type of vehicle, using a commercial power supply. Generally, the commercial power supply has high power generation efficiency.\nAs one type of the hybrid vehicles, a vehicle that permits an on-board electric storage device to be charged using electric power from a power supply located outside the vehicle (which may also be simply called “external power supply”), like an electric vehicle, is known. The charging of the on-board electric storage device with electric power from the external power supply may also be simply called “external charging”. Typically, the vehicle as described above is known as so-called “plug-in hybrid vehicle”. For example, a charge port or socket provided on the plug-in hybrid vehicle is connected to an outlet of a power supply provided in a house, or the like, via a charging cable. With this arrangement, the on-board electric storage device can be charged with electric power from a standard home power supply. By using the external power supply, it is expected to enhance the fuel consumption efficiency of the hybrid vehicle.\nJapanese Patent Application Publication No. 2007-178215 (JP 2007-178215 A) discloses a power supply system capable of charging a secondary battery. The power supply system restricts charging current in view of the temperature dependence of the internal resistance of the secondary battery. With the charging current thus restricted, the charging status (which may also be called “state of charge” or “SOC”) of the secondary battery is estimated with improved accuracy, as described in JP 2007-178215 A.\nA vehicle capable of external charging assures a long running distance if the electric storage device is in a sufficiently charged state when driving of the vehicle is started next time. In the meantime, if the electric storage device is kept in a high SOC condition for a long period of time, degradation of the electric storage device may be accelerated. Therefore, a vehicle having a timer charging function has been proposed. The user of the vehicle sets a scheduled time at which charging is to be completed (or a scheduled time at which driving of the vehicle is to be started). The charging of the electric storage device is controlled so as to maximize the SOC immediately before the scheduled time set by the user.\nIn the vehicle having the timer charging function as described above, if the estimated charging time of the electric storage device is not accurate, a shortage of charge may occur at the scheduled charging completion time, or charging may be completed earlier than scheduled. In the latter ease, the electric storage device may be kept in a high SOC condition for a long period of time, which may result in degradation of the electric storage device.\nIn the meantime, charging power that can be received by the electric storage device may be restricted by the temperature and SOC of the electric storage device. Therefore, it is important to take account of charging characteristics of the electric storage device, for accurate estimation of the charging time."}
{"text": "Modern electronic design is typically performed with computer aided design (CAD) tools or electronic design automation (EDA) systems. To design a circuit, such as an integrated circuit (IC), a designer first creates high level behavior descriptions of the IC device using a high-level hardware design language (HDL). Common examples of HDLs include Verilog and VHDL. An EDA system typically receives the high level behavioral descriptions of the IC device and translates this high-level design language into netlists of various levels of abstraction. Essentially, the process to implement an electronic device begins with functional design and verification (e.g., using RTL), and then proceeds to physical design of a layout and verification.\nCircuit designers and verification engineers use different methods to verify circuit designs. One common method of verification is the use of simulation. Simulation dynamically verifies a design by monitoring behaviors of the design with respect to test stimuli. For many types of designs, simulation can and should be performed during the design process to ensure that the ultimate goals are achievable and will be realized by the finished product.\nFor example, SPICE and FASTSPICE are common types of simulators that perform simulation to verify the operation of an electronic design. With these types of tools, the electronic design is converted into a system of equation(s), which is then solved for a given set of inputs to check the state of specific portions of the circuit at given points in time.\nA significant portion of the data pertaining to the simulation activities may be waveform data. In many cases, the waveform is the output data that results from performing simulation on analog or analog/mixed signal designs. For example, a given circuit design may receive a set of stimulus data, where simulation of a model of that circuit design is performed using the stimulus to identify the resulting waveform pattern from expected operation of that design. In other cases, the waveform data is the input data that is fed into a circuit design for simulation.\nThe EDA tool may employ either a waveform viewing tool or a waveform editing tool to visually display waveform data onto the user interface of a user display apparatus. In many cases, these tools are operated in an interactive manner, where a user dynamically manipulates the tool to display the waveform data in real-time with a desired visual representation. For example, the user may manipulate the waveform tool to zoom into a specific portion of the waveform or to scroll across different portions of the waveform. This interactive manipulation permits the user to analyze the waveforms and to seek out anomalies that may indicate possible errors in the circuit design, by scanning through the data and zooming in at specific portions of the waveforms.\nUnder the covers, numerous computational operations may be performed to comply with the user's manipulation of the waveform tool. Such operations may involve loading the waveform data into memory, retrieving the specific data to be operated upon, performing graphical operations upon that data to generate display data as required by the user commands, and generation of the interface for display. With interactive tools, it is desirable for these operations to be performed fast enough such that the user does not experience undue delays while rendering the waveform for display on the display screen.\nThe problem is that modern EDA tools are now generating massive quantities of waveform data, e.g., involving tens of millions of datapoints on thousands of waveforms. In many cases, many or all of the waveforms need to be displayed simultaneously onto a computer display screen. Given this large volume of waveform data, it is becoming very difficult to maintain interactive performance from conventional waveform viewing tools. These performance problems occur because even simple calculations become very difficult to process when there is both a large amount of data and a need to quickly process that data for display on an interactive tool.\nTherefore, there is a need for an improved approach to handle waveform data that that can adequately address these problems with the conventional tools."}
{"text": "1. Field of the Invention\nThe present invention relates to shock-absorbing structures of bicycles and more particularly, to a shock-absorbing seat stay for a bicycle.\n2. Description of the Related Art\nWhen riding a bicycle on a bumpy road, a bicyclist may feel uncomfortable due to vibration. To ease discomfort of the bicyclist, many bicycles may be equipped at the bicycle frame thereof with a suspension for buffering the vibration. However, the conventional suspension for being installed with the frame of a bicycle usually comprises metal springs, pneumatic device or hydraulic device for achieving the purpose of absorbing vibration, resulting in that the aforesaid conventional suspension is complicated in structure and heavy in weight, and even a part of stepping energy from the bicyclist may be absorbed by the conventional suspension to undesiredly deteriorate the transmission efficiency of energy. Therefore, for a road bicycle built for traveling at high speed with a light-weight specification, the aforesaid conventional suspension composed of metal spring, pneumatic device or hydraulic device is not an appropriate choice. On the other hand, when a bicyclist rides a road bicycle without suspension, the bicyclist may easily and quickly feel tired. Under this circumstance, the bicyclist may have to face a dilemma that he/she can only choose one from a bicycle having high energy transmission efficiency with light weight and a bicycle having suspension that can provide riding comfort."}
{"text": "Data processing systems may use virtual addresses to indicate storage locations while the processing system uses physical addresses that represent actual locations on the silicon to store the data in. Virtual addresses may be used to reduce the number of bits required to identify an address location or to allow several processes to access a restricted memory space.\nA complete set of current mappings for the virtual to physical addresses are stored in memory, however in order to decrease access time to these mappings, recently used mappings are stored in caches that can be accessed more quickly by the processor. There may be an L1 cache that is fast to access and acts as a micro TLB and stores a small subset of recently used mappings, and a slower L2 cache that is the macro TLB and stores a set of currently used mappings while a full set of page tables of mappings are stored in memory.\nThe mapping of virtual to physical address space is done in blocks and these may vary in size, thus, these are for example blocks of 1 Gbyte, 2 Mbyte or 4 Kbyte. The number of bits that need to be mapped for a translation depends on the size of the block that an address is located in. If, for example, an address is in a 2 Mbyte block, then only the higher n to 21 bits need to be found from the translation tables, while if it is in the 4 Kbyte block then the n to 12 bits need to be found.\nThe mappings and the final page sizes are therefore stored in tables in the memory, with the first n to 31 bits that represent 1 Gbyte blocks, being stored in a first table, the next 30 to 21 bits being stored in a next table and so on. A page table walk in memory that is used to retrieve the mappings is performed in steps or walks and where the address is in a larger block only the first step(s) needs to be performed.\nWhen a mapping has been found following a page table walk it is stored in the L2 macro TLB and in the L1 micro TLB.\nAs can be appreciated retrieving these mappings from memory is expensive in both time and power.\nIn order to increase efficiency in some data processing systems, the RAM used as the level 2 cache for the macro TLB is also used as a walk cache and as a prefetch buffer, such that in addition to retrieving the requested mapping, a mapping for the subsequent virtual address is retrieved during the same page table walk. This can then be written into the prefetch buffer, after the originally requested mapping has been written into the macro TLB.\nThe walk cache will store intermediate translations following some of the page table walks, the stored values being an input to a final page table walk required to finish the translation.\nIn some systems there may be double virtualisation. This occurs where there are several guest OSs and each references a large amount of memory. The memory of the data processing system is not sufficiently large from them to all to reference different distinct portions of memory. Instead they share some memory area but this is not visible to them and is managed by the hypervisor with each guest OS using intermediate physical addresses IPA, that the hypervisor then maps to real physical addresses. In this way the hypervisor controls the memory use and allows each guest OS to believe it has sight of a large dedicated memory space. In this case each guest OS manages the page tables for the VA to IPA translations and the hypervisor manages tables of translations from IPA to PA. In these cases an IPA2PA cache may be provided in the level 2 RAM storing intermediate steps in the IPA to PA translations."}
{"text": "The present invention relates to protection of nuclear fuel assemblies and, more particularly, to a method of and an arrangement for protecting a fuel assembly such as to avoid damaging and scratching of the fuel rod surfaces during assembly and transportation of the fuel assembly.\nDifferent types of fuel assemblies are used in different types of nuclear reactors such as boiling water reactors (referred to as \"BWR\", hereinunder), pressurized water reactors (PWR) and so forth. Generally, however, a fuel assembly is an elongated member in which a plurality of elongated fuel rods are arranged and held in parallel with one another by spacer assemblies.\nFIGS. 1 to 3 show, by way of example, a fuel assembly for use in a BWR having a square cross-section with fuel rods arranged in a square pattern. Known also is a fuel assembly having a circular cross-section in which fuel rods are arranged along concentric circles. This type of fuel assembly is suited for use in heavy-water moderated reactors and gas-cooled reactors.\nThe spacer assembly for holding the fuel rods in parallel can have various forms such as rectangular grids and circular grids. The fuel rods are held in the cells of the spacer assembly by means of holders such as a spring-clip type holder and lantern-ring type holder.\nThe fuel assembly shown in FIG. 1 will be explained in detail. This fuel assembly has a channel box 1 accomodating a plurality of fuel rods 4 and water rods 5 which are supported at their upper and lower ends by an upper tie plate 2 and a lower tie plate 3. A plurality of spacer assemblies 6 are disposed in the channel box 1 in the axial direction of the latter, such as to hold the fuel rods 4 in parallel and to suppress any lateral vibration of the same, while preserving passages of the coolant.\nAs will be seen from FIGS. 2 and 3, the spacer 6 has a square frame 7 in which are disposed bars 8, 9 and dividers 10, 11 in a grid-like form thus defining a plurality of cells 12. A lantern spring 13 is provided at each point of intersection of the bars 8 and 9. Each lantern spring 13 has leaf springs which project outwardly from all sides thereof. The leaf springs 13a resiliently press the fuel rod 4 in each cell 12 against S bends 10a, 11a formed on the dividers 10, 11 and dimples 14 projecting inwardly from the walls of the frame 7.\nFor the assembly of the fuel assembly, the lower tie plate 3 and the spacer assemblies 6 are fixed at suitable positions and the fuel rods 4 and the water rods 5 are inserted into the cells 12 of the spacer assemblies 6, such that the lower ends of the rods 4, 6 are fixed in the lower tie plate 3. Then, the upper tie plate 2 is mounted on the upper end of the frame 7 thus fixing the fuel rods 4 and the water rods 5.\nDuring the assembly of the fuel assembly, the surfaces of the fuel rods tend to be damaged because they slide on the surfaces of the leaf springs 13a of the lantern spring 13 and the S-bends 10a and 11a. The same problem is encountered not only by the spacer assembly shown in FIG. 2 but also by various other types of spacer assemblies.\nHitherto, various methods have been proposed to protect the surfaces of the fuel rods. For instance, it has been proposed to tie the leaf springs 13a and S-bends 10a, 11a by cords against their resiliency thus preserving a large area in each cell 12 for receiving the fuel rod. It has been also proposed to keep the fuel rod below the ice point such as to form a coating layer of ice or frost on the fuel rod surface. According to still another method, the fuel rods are beforehand painted with a lacquer and, after the assembly, the lacquer is washed away.\nThe protection method relying upon the ice or frost coating layer necessitates pre-cooling of the fuel rod down below the ice point. In addition, this method is not adoptable for the fuel rod protection during withdrawal which often becomes necessary when, for example, the defective fuel rod is renewed, because the formation of the ice or frost coating layer on the fuel rod situated in the fuel assembly is almost impossible. The protection method using the cords is quite troublesome and takes a long time. The use of a lubricant may impair the performance of the fuel assembly due to deposition of impurities on the fuel rod. In order to completely remove the lubricant, it is necessary to wash it away at a cost of time and labour.\nThus, the conventional fuel rod protection methods are generally laborious and hinders the assembling.\nThe damaging of the fuel rods in the spacer assembly takes place not only during the assembly but also during transportation. Namely, during transportation of the completed fuel assembly, the fuel rods which are held in pressure contact with the S-bends 10a and 11a by means of the resilient pressure members such as the leaf springs are liable to be fretted due to vibration of the fuel assembly. To avoid this problem, it has been proposed to use flexible plastic liner 21 of a length almost equal to that of the fuel rod. Namely, as shown in FIG. 4, a liner 21 is placed at each side of the fuel rod in each cell, so as to suppress the vibration of the fuel rod and to prevent fretting of the fuel rod surface by the spacer assembly. However, the fuel assemblies 4a adjacent the walls of the frame 7 are likely to be damaged due to vibration, because these rods are pressed onto the dimples 14 by the leaf springs 13.\nJapanese Patent Laid-Open No. 4609/1976 discloses a rod-like jig having a plurality of tabs designed for pressing the springs 13. On the other hand, Japanese Patent Laid-Open No. 36282/1977 and Japanese Patent Laid-Open No. 16089/1979 disclose jigs which are designed to be inserted into cells 12 such as to press the springs 13.\nThese known jigs are intended to retract the springs 13 in advance to the insertion of the fuel rods, thereby smoothing the movement of the fuel rods during insertion. However, these jigs cannot prevent the damaging of the fuel rods perfectly because, although the springs are retracted beforehand to become inactive, the fuel rods slide on the S-bends 10a, 11a and the dimple 14 when they pass through the cells 12. In some cases, the spring 13 is adversely affected by an excessive depressing by the jigs.\nJapanese Patent Publication No. 11244/1978 discloses a fuel rod protection method in which a stainless steel protecting tube of about 0.1 mm thick is fitted around the fuel rod, and the fuel rod is inserted into the cell 12 together with this protecting tube. However, insertion and withdrawal of the fuel rod into and out of such a thin protecting pipe are troublesome. In addition, the fabrication and handling of such thin-walled long tube are quite difficult."}
{"text": "A standard wheelchair seat consists of a flexible sling material. The material has adequate strength to support the weight of the user. However, the basic sling shape creates dangerous high pressure areas to the user's skin. The prolonged high pressure can cause pressure sores leading to potentially fatal Pressure Ulcers. The typical method to improve the pressure distribution and support of a person in a wheelchair is to add an additional seat cushion system on top of the wheelchair seat sling. There are many different types of cushioning and support systems. They may use air bladders, fluid filled bags, different density foams, contoured foams or a combination of several technologies. Examples of these methods are shown in U.S. Pat. Nos. 5,687,436; 5,282,286; 4,726,624; and 5,839,140.\nAnother method is shown in U.S. Pat. Nos. 6,755,477 and 6,264,279 in which multiple layers of flexible non-uniform material such as Lycra or Polartec are specified. The non-uniform materials are used in combination with flexible and/or stiffened materials between the layers to give the seat contours to improve the pressure distribution. Another method is shown in U.S. Pat. Nos. 5,857,749 and 5,984,418 in which the flexible material is replaced with a rigid base on which a cushion is attached. All the known methods come at the price of added cost, weight and complexity to the seating system."}
{"text": "1. Field\nApparatuses and methods consistent with exemplary embodiments relate to controlling travel of a vehicle in which a driving mode of the vehicle is determined based on the speed of the vehicle and a user-selected driving mode, and more particularly to controlling travel of the vehicle when a failure occurs in a power supply that supplies the vehicle with power according to the determined driving mode, the power supply is replaced with another power supply for stable driving.\n2. Description of the Related Art\nJapanese Patent Application Publication No. 2003-032807 discloses that a vehicle is switched by a hybrid vehicle control device between an electric vehicle (EV) mode in which the vehicle travels by using a battery alone, an engine mode in which the vehicle travels by using an engine alone, and a hybrid electric vehicle (HEV) mode in which the vehicle travels by using the battery and the engine together according to stored road data.\nSuch related art technology merely enables a switch between the three driving modes based on road data, and does not reflect the output or battery/engine state of a vehicle. In addition, a means and method for storing data about a driving region are necessary for the related art. In other words, the related art will reach to the limit of storage capacity and involve additional cost for the storage means."}
{"text": "1. Technical Field\nThe present disclosure relates to a portable electronic device, and more particularly to a portable electronic device capable of hiding a flexible printed circuit thereof.\n2. Description of Related Art\nIn a conventional portable electronic device such as a slider computer, the screen is slidably disposed on the keyboard. As the slider computer in use, the screen is slide to one side of the keyboard exposing a keyboard thereof and upwardly being lifted on the keyboard. However, in sliding the screen of the slider computer to stand on the keyboard, an electric connecting strip (e.g., a flexible printed circuit) for connecting the screen and the keyboard or host would be exposed outside the slider computer, so that the electric connecting strip may have higher possibility of being damaged.\nIn this regard, how to develop a portable electronic device for solving the mentioned disadvantages is an urgent issue in the art."}
{"text": "Internal combustion engines may include water injection systems that inject water from a storage tank into a plurality of locations, including an intake manifold, upstream of engine cylinders, or directly into engine cylinders. Injecting water into the engine intake air may increase fuel economy and engine performance, as well as decrease engine emissions. When water is injected into the engine intake or cylinders, heat is transferred from the intake air and/or engine components to the water. This heat transfer leads to evaporation, which results in cooling. Injecting water into the intake air (e.g., in the intake manifold, ports, etc.) lowers both the intake air temperature and a temperature of combustion at the engine cylinders. By cooling the intake air charge, a knock tendency may be decreased without enriching the combustion air-fuel ratio. This may also allow for a higher compression ratio, advanced ignition timing, and decreased exhaust temperature. As a result, fuel efficiency is increased. Additionally, greater volumetric efficiency may lead to increased torque. Furthermore, lowered combustion temperature with water injection may reduce NOx, while a more efficient fuel mixture may reduce carbon monoxide and hydrocarbon emissions.\nAs explained above, water may be injected into different locations, including the intake manifold, intake ports of engine cylinders, or directly into engine cylinders. The inventors have recognized that water injection at different locations may provide different benefits depending on engine operating conditions. For example, manifold injection, direct injection, and port injection toward the manifold may provide increased cooling of the charge air, thereby reducing knock tendency. Additionally, port injection onto the intake valves may increase a dilution effect of the injected water over manifold injection, thereby reducing engine pumping losses. However, an actual amount of water injected at an engine may be different than a commanded amount. For example, water injection errors may arise due to water injector fouling or wear, due to incomplete vaporization of the injected water, which may result from impurities in the water being injected, variation of the injected water pH level from an expected level, a CAC outlet temperature being different than expected, etc. As a result, a desired water injection benefit may not be provided to the engine.\nIn one example, the issues described above may be addressed by a method for an engine comprising, responsive to water injection into an intake manifold via a port water injector, adjusting engine operation based on a change in exhaust dilution; and responsive to water injection into the intake manifold via a manifold water injector, adjusting engine operation based on a change in intake dilution. In this way, water injection errors may be diagnosed differently based on the location of the water injection as well as based on the intended benefit of the water injection.\nAs an example, an engine may be configured with multiple water injectors such as a manifold water injector, a port water injector, and a direct water injector. An engine controller may determine a total amount of water to inject and a location to inject the water based on engine operating conditions. For example, at low-mid engine loads, engine combustion stability may be improved by increasing the charge dilution via injection of water into the intake manifold. Additionally, a portion of the total water to be injected may be port injected onto a hot surface of a closed intake valve to meet the dilution demand. Since charge dilution via water injection affects the amount of water in the engine, the actual amount of water injected may be detected via an oxygen sensor. Specifically, a change in the output of an intake oxygen sensor may be used to infer the actual amount of water injected into the intake manifold, while a change in the output of an exhaust oxygen sensor may be used to infer the actual amount of water injected into the intake port. By comparing the actual water injection amount to the commanded amount, the controller may learn a water injection error and compensate accordingly. For example, during a subsequent water injection, a duty cycle of the manifold and/or port injector may be adjusted to compensate for the error. As another example, EGR flow to the engine may be adjusted to compensate for any dilution deficit. In comparison, during mid to high engine load conditions, knock control may be provided by using manifold water injection to provide charge cooling. During such conditions, the water injection error may be learned as a function of a change in the manifold temperature (via a manifold temperature sensor), and compensated for via adjustments to the duty cycle of the manifold injector, via direct injection of water into a cylinder, and/or spark timing adjustments.\nIn this way, by selecting a water injection sensing mode based on the engine load and the selected water injector(s), water injection errors may be more accurately measured and compensated for. The technical effect of using a different set of sensors to measure water injection used for charge dilution (at low-mid engine loads) versus water injection used for charge cooling (at mid-high engine loads), water injection may be better controlled. By providing a desired dilution to the engine, pumping losses may be decreased and combustion stability may be improved. By providing a desired charge cooling to the engine, knock related issues may be reduced. Overall, water injection benefits may be extended over a wider range of engine operation, improving engine performance.\nIt should be understood that the summary above is provided to introduce in simplified form a selection of concepts that are further described in the detailed description. It is not meant to identify key or essential features of the claimed subject matter, the scope of which is defined uniquely by the claims that follow the detailed description. Furthermore, the claimed subject matter is not limited to implementations that solve any disadvantages noted above or in any part of this disclosure."}
{"text": "1. Technical Field\nThe present invention relates to a system and method for generating voice messages from pre-recorded audio segments. Individual pre-recorded audio segments are dynamically cross-faded during playback of audio messages synthesized from a plurality of such audio segments. Dynamically cross fading the sound segments speeds up the playback time of the voice messages and provides more natural sounding voice messages.\n2. Related Art\nIn many automated systems in which a human user must interact with a computer through an audio interface, the automated system communicates with the user via audible voice prompts. Such voice prompts are typically created by chaining together multiple pre-recorded audio segments. For example, an automated telephone information service communicates telephone numbers to a requester by sequentially playing audio files representing each digit of a telephone number. Similarly, most voice recognition systems require feedback to the user. Typically such feedback is in the form of audio voice prompts requesting certain actions or information from the user. The content of such voice prompts is often context specific, so the system must select the appropriate pre-recorded sound segments to join together “on the fly” from among multiple pre-recorded sound segments in order to create a message that delivers the desired information, or asks for the appropriate response from the user.\nTypically, such voice messages are created by simply concatenating a plurality of individual audio files and sequentially playing them back over a loudspeaker. Often the individual sound files will have short silent lead-in and exit times on each side of the actual sound portion of the audio file. When multiple files are played back, these short silent portions result in a message that sounds choppy and unnatural and takes much longer to play back than would otherwise be the case if the message were spoken by a live person. Such a delay and unnatural delivery can be a source of frustration for users of systems that employ synthesized voice prompts. Often a user must wait for the entire message prompt to be played back before being allowed to proceed with the next step in an automated transaction. The undue delay from excessively long voice prompts can be tedious and annoying.\nFIG. 1 shows three typical pre-recorded sound files 10, 12, 14 representing the spoken words “six”, “seven” and “eight” respectively. Each of the sound files 10, 12, 14 includes a lead in time 16, and an exit time 18 where no actual sound is produced. FIG. 2 shows the same files concatenated together to form the sequence “six-seven-eight”. As can be seen, there are relatively long periods of silence 22 between the combined sound files. No information is provided to the user during these intervals. These delays unduly lengthen the playback time and detract from the sound quality of the resulting voice message.\nFIG. 3 depicts a method of clipping audio files so that only those portions of the audio file in which sounds are actually voiced are included in synthesized voice messages. Unclipped sound segment 24 has silent lead in and exit portions 16, 18. In the clipped sound segment 26 the silent portions have been removed. As can be seen, the duration of the original unclipped file 24 is significantly longer than the clipped file 26. FIG. 4 again shows the numbers sequence “six-seven-eight” 28 but this time composed of sound segments that have been clipped in the manner shown in FIG. 3. Clipping the sound segments in this manner shortens the playback time of the three digit sequence from 2.12 seconds to 1.93 seconds. Although this represents a significant improvement over the unclipped sequence, the result is still longer and sounds choppier than when the sequence is spoken by a natural person. Thus, there remains a need for improved systems and methods for synthesizing voice messages and prompts from individual sound files. Such improved systems and methods must produce shorter more natural sounding voice messages and prompts."}
{"text": "Parametric stereo and multichannel coding methods are known to be scalable and efficient in terms of listening quality, which makes them particularly attractive in low bitrate applications. In cases where the bitrate limitations are of a transitory nature (e.g., network jitter, load variations), however, the full benefit of the available network resources may be obtained through the use of an adaptive distribution format, wherein a relatively higher bitrate is used during normal conditions and a lower bitrate when the network functions poorly. Existing adaptive distribution formats and the associated (de)coding techniques may be improved from the point of view of their bandwidth efficiency, computational efficiency, error resilience, algorithmic delay and further, in audiovisual media distribution, as to how noticeable a bitrate switching event is to a person enjoying the decoded media.\nAll the figures are schematic and generally only show parts which are necessary in order to elucidate the invention, whereas other parts may be omitted or merely suggested. Unless otherwise indicated, like reference numerals refer to like parts in different figures."}
{"text": "Using devices, such as active and passive air fresheners, including electrical plug-in air fresheners, candles, aerosol diffusers, sprays, membrane air fresheners and the like, to deliver nice and refreshing fragrances to air within a room is well known. There are numerous devices of this kind available on the market place, with a number of fragrances available as well. The targeted consumer benefits include, for example, pleasantness of the neighbor odor, hedonics, malodor counteracting, and insect repellency.\nForest bathing, which is a popular exercise in Asian countries, is a means of body sculpting, which includes taking walks and doing aerobic exercises in a forest. It is believed that forest bathing trips have a positive impact on human beings, such as stress and anxiety relief, calming and reinvigorating actions. It is recognized that these benefits are related to the presence of phytoncides or wood oil essence components in the atmosphere of forests. In particular, a correlation between the presence of phytoncides and the activity of the human natural killer cells has been documented. See for example Q. Li at al. International Journal of Immunopathology and Pharmacology 22 (2009) pages 951 to 959 and references cited therein. The enhancement of the natural defense of the organism induced by the exposure of individuals to the atmosphere of natural forests, especially those populated by resinous trees or Eucalyptol trees is also well known. These phytoncides are volatile molecules of the terpene and sesquiterpene families, which are produced by the wood itself, whereas such a production is especially important in resinous trees, such as pine trees and cypress trees.\nHowever, people are busy with work and often do not have the time to travel to a forest to enjoy the air environment in a forest, especially, the inhabitants in the cities. Accordingly, there remains a need to provide an air freshener device, which generates and delivers the health and therapeutic benefits of forest air to the inside environment."}
{"text": "In U.S. Pat. No. 4,232,184 which was issued to me on Nov. 04, 1980, I show a cable adapter, commonly referred to in the trade as a boot, for use with a cable closure at each of its ends in aerial telephone plant. The cable closure protects a cable splice. Likewise, some aerial cable terminals may use a boot at each end. The terminal provides a location for terminating a loop.\nThe boot which I have described in the aforesaid patent has a plurality of tapered nozzles for adapting to the size of the cable. Initially, the boot was made of soft material. But rodents could eat through the soft material, exposing the cable to elements and resulting in potential maintenance problems. An improved boot was made from a hard substance. This, however, made it difficult for a craftperson to work with the apparatus."}
{"text": "Alkyl lead compounds, which have been conventionally used as anti-knock additives to motor fuels, exhibit undesirable properties such as the toxicity of the lead and their poisoning effect on catalytic converters. As a result, a need has arisen worldwide for other high octane compounds to be blended into the pool of fuel components available from refineries and synfuel plants.\nBecause of certain limits to the extent to which oil refineries and synfuel plants can increase the octane number of available fuel pool components, the demand for isoalkenes, especially in the C.sub.4 to C.sub.5 range, has risen considerably in recent years. One of the products which can be synthesized from isobutene is methyl tert-butyl ether (MTBE), which is a valuable octane boosting fuel additive. As a consequence, there is an increased demand for isobutene and for isomerization processes for the conversion of linear butenes to isobutene. Similarly, linear pentenes can be converted to isopentenes which could then be used to produce tert-amyl methyl ether (TAME).\nlsomerization processes can be directed towards either skeletal isomerization or double bond isomerization. Skeletal isomerization is concerned with the reorientation of the molecular structure of a hydrocarbon, so as to increase the number of side chains. Double bond isomerization is concerned with the relocation of a double bond between carbon atoms forming part of a chain, and is not of great interest in the context of increasing octane values.\nA large number of isomerization processes has been described in the literature, which has been summarized by Butler and Nicolaides in Catalysis Today, Vol. 18, 1993, 443 -471. Some of these processes rely on the use of metal oxides or crystalline aluminosilicates, some of them halogenated and others not, some of them containing metals and others not, some of them requiring a diluent in the isomerization reaction and others not. Most of the zeolite and halogenated alumina catalysts for the isomerization of linear butenes have a short life due to the formation of coke during the isomerization reaction, resulting from the occurrence of parasitic reactions such as cracking, polymerization and oligomerization, and a lack of stability frequently resulting in a rapid decrease of the conversion rate.\nRSA Patent Application No 92/1318 and European Patent Application No 92200516.0 disclose a process for the conversion of a feedstock comprising linear olefins into a product enriched in branched olefins, which process comprises contacting the feedstock with a tectometallosilicate having a ferrierite crystal structure at a temperature between 150.degree. C. and 450.degree. C., and olefin partial pressure of more than 0.5 bar and a total pressure of between 0.5 and 25 bar. As will be appreciated by those skilled in the art, a tectometallosilicate which has a ferrierite crystal structure, is a fully or at least substantially crystalline material.\nThe use of the ferrierite zeolite, as a catalyst for the skeletal isomerization reaction of linear buteries, at an olefin partial pressure of Jess than or equal to 0.5 bar is described in European Patent Application 9320 1677.7. As expected, the lower partial pressures result in higher yields of the branched isomer.\nEuropean Patent Application No 92305090.0 discloses a process for structurally isomerizing a linear olefin of at least 4 carbon atoms to its corresponding methyl branched iso-olefin, which process comprises contacting a hydrocarbon feedstream containing at least one said linear olefin with an isomerizing catalyst at a temperature of from 340.degree. C. to 650.degree. C., said isomerizing catalyst comprising at least one zeolite with one or more one-dimensional pore structure having a pore size small enough to retard by-product and coke formation within the pore structure and large enough to permit entry of the linear olefin and allow formation and diffusion of the methyl branched iso-olefin. Likewise, it will be appreciated that a zeolite is considered by those skilled in the art, as a fully crystalline or at least substantially crystalline material. Examples of zeolites which can be used according to the aforementioned publication, include ferricrite, FU-9, NU-23, NU-10, ZSM-12, ZSM-22 etc, and molecular sieves such as SAPO-11, SAPO-31, SAPO-41, etc."}
{"text": "The present invention relates generally to apparatuses and methods for metering and conveying dry powder or granular materials to a blender, and more particularly to an apparatus and method for accurately metering and conveying dry powder or granular materials to a blender in a substantially closed system.\nIn conventional oil field operations, dry powder and granular materials are transported in sacks to the well location. The sacks are then manually moved to a metering feeder mounted on a fluid blending device. The metering feeder has a hopper to collect the dry powder or granular material and to serve as a reservoir. The hopper is kept full by opening sacks of the dry powder or granular material and dumping them into the hopper. The amount of dry powder or granular material must be kept at a constant level in the hopper so as to maximize the accuracy and repeatability of the metering feeder. The accuracy and repeatability of the metering feeder is dependent on having a constant material bulk density at the screw auger.\nFluctuating levels of powder or granular material in the hopper can alter the bulk density of the material, and thus negatively affect the desired quality of the oil field service fluid. However, it is difficult to maintain a constant feed of material to the hopper when the hopper is being filled manually. Thus, the drawback of such systems is that there is a likelihood the hopper will be either overfilled in which case dry powder or granular material is spilled, or underfilled in which case the bulk density will be altered, which in turn will reduce the quality of the oil field service fluid. In the case where the dry powder or granular material is spilled, material is lost and clean up is necessary, which translates into lost revenue. This method also generates a lot of dust since the hopper is open to the outside air.\nAlternate methods have been proposed, which also have drawbacks of their own. One such method is to transport the dry powder or granular material to the well location in a large bulk tank and to convey the material pneumatically or mechanically to the metering feeder on the fluid blending device. This method, however, requires the conveying device to be controlled either automatically with level sensors and control loops, or manually in order to keep the dry powder or granular material at a constant level in the metering hopper for maximum accuracy and repeatability. Vehicle roading and equipment vibrations tend to compact the powder or granular material and cause significant variations in the bulk density of the material which affects the accuracy and repeatability of the metering feeder. This process is also labor intensive, susceptible to sensor failures and likely to generate dust."}
{"text": "The present invention relates to a ceramic microwave filter having improved edge steepness, which has at least two resonators in a one-piece ceramic body and coupling structures associated with the resonators, for capacitively coupling the ceramic filter to an external radio-frequency circuit.\nDE 197 41 147 A1 describes a ceramic microwave filter comprising a ceramic body in which two mutually parallel inner conductor bores with a round cross section are constructed as resonators. The ceramic body itself comprises two sintered green compacts pressed together, and is therefore seen to that extent as being one piece although the two green compacts form different dielectrics. The ceramic body also has coupling structures which are formed by electrical isolation of the metallic coating which is otherwise applied over the whole surface to the ceramic body. This isolation can be undertaken mechanically or by an etching technique or photo-resist technique.\nFurthermore, DE 197 41 177 A1 describes a ceramic microwave filter which is formed from a standard ceramic body having coupling structures constructed on its surface and two inner conductor bores with a circular cross section as resonators, these two inner conductor bores having steps which serve to steepen the flanks of the forward characteristic of the ceramic microwave filter. The ceramic body of this ceramic microwave filter is also metallized on all sides with the inclusion of the inner surfaces of the inner conductor bores, the coupling structures being electrically isolated from this metal coating. Such ceramic microwave filters are preferably used as mobile radio filters in mobile radio telephones (mobiles).\nMobile radio sets and, specifically, mobile radio telephones are intended to have as high an output power as possible in order to be able to ensure a maximum output power in mobile radio sets, transmitting filters should be distinguished by insertion losses that are as low as possible, on the one hand, and transmission noise in a receiving band should be minimized by transmitting filters with a high selectivity, on the other hand. Thus, overall there is a need for a transmitting filter which has a high selectivity in conjunction with low insertion losses.\nIn the case of ceramic microwave filters, there is a fundamental relationship between insertion loss or the quality of the filter and the volume of the ceramic body of the filter. The larger the ceramic body made from a highly dielectric ceramic, the lower the insertion loss, since the quality of the line resonators coupled in the ceramic body is proportional to the volume of the ceramic body. This means that a reduction in the insertion loss can be achieved only via an increase in the volume of the ceramic body. However, this is opposed by the requirement that the ceramic microwave filters should be fashioned as small as possible for the purposes of mobile radio, in order likewise to be able to fashion the size of a mobile radio telephone to be as small as possible. It is to be noted that the selectivity can be H achieved only by increasing the pole number or by a complicated coupling of the resonators of the ceramic microwave filters with coils and capacitors to form a band-stop filter.\nIt is accordingly an object of the invention to provide a ceramic microwave filter having improved edge steepness which is distinguished by a lower insertion loss at the band edge, and thus by a correspondingly high edge steepness.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, a ceramic microwave filter having a characteristic curve with an improved edge steepness. The ceramic microwave filter includes a one-piece ceramic body with at least two resonators and with coupling structures associated with the resonators for capacitively coupling the ceramic filter to an external RF circuit. At least one counter-oscillator is configured in the ceramic body approximately in parallel with the resonators.\nIf appropriate, it is also possible to provide two or more counter-oscillators, it being preferred for each resonator to be assigned a counter-oscillator. The resonators themselves can be rectangular in cross section, since given such a cross section, the unloaded quality, and thus the insertion loss of the ceramic microwave filter, can be further improved.\nThe ceramic body is completely metallized except for its raised end face, the coupling structures being isolated in the usual way from this metallic coating.\nThe two resonators which, as explained, preferably have a rectangular cross section, are excited capacitively via the coupling structures. The counter-oscillators are weakly excited, in order to render the insertion loss of the frequency at which the xe2x80x9ccounter-notchxe2x80x9d forms in the characteristic curve as highly selective as possible. Specifically, at the same time the insertion loss at the band edge is then low, with the result that the edge steepness is correspondingly high.\nThe edge steepness and selectivity can be further increased by use of a second counter-oscillator when the xe2x80x9ccounter-notchxe2x80x9d is selected to be approximately the same for both resonators, Another possibility consists in providing the counter-oscillator such that both filter edges of the ceramic microwave filter are steepened.\nIt is also possible, in principle, to use more than two resonators, although the insertion loss then rises slightly. However, a further improvement in the edge steepness can also be achieved here by using appropriate counter-oscillators.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in a ceramic microwave filter having improved edge steepness, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to articulating fencing systems and, more particularly, to a radially articulating fencing system.\n2. Description of the Related Art\nFarmers and ranchers use fences, walls, hedgerows and other boundary markers to demarcate territory, and to control ingress and egress of persons and animals from and to the boundaries territory. With the advent of large cattle and sheep ranches, particularly in the Western United States, fences are used primarily to restrain domestic livestock from leaving a fenced in area. These western style fences are typically built from four or five horizontally spaced strands of barbed wire, mounted on wooden fence posts that are driven into the ground. A single fence, often stretching for miles, is used to delineate a xe2x80x9crangexe2x80x9d or xe2x80x9cpasturexe2x80x9d where herds of cattle, or other grazing animals, are confined.\nMulti-purpose land use methods have brought crop cultivation and the raising of livestock together on the same land. These multi-use methods involve dividing a tract of land into several parcels and rotating the function of the parcels between crop cultivation and livestock feeding. These methods typically involve concentrating dense populations of livestock, such as feeder cattle and dairy cattle, on one parcel while crops, such as grains and seeds, are being cultivated on the others. Since cattle fed in this manner are in high density, barbed wire fences, and fences that have electrified strands of wire, are needed to confine the cattle in the parcel. Multi-purpose land use methods have gained widespread acceptance in the agriculture business because they generate increased yields of table meat from cattle by providing abundant food supplies for fattening cattle prior to slaughter.\nWith the requirement for increased agriculture production, especially in arid climes, sophisticated irrigation systems and methods have become a necessity to increase acreage yields. Today, these systems are usually automated, using computers, and can cover extended acreage without being manned. Water pressure and electricity are commonly used to provide the energy to move these systems. The irrigation systems which have mechanized means of movement such as drive wheels mounted under spray irrigation carriages. A pivotal irrigation system, for example, is anchored at a center point and rotates about that point on large cleated wheels, mounted under a spray irrigation carriage, to sweep out an irrigation circle which can be a mile or more in diameter.\nThe close proximity of cattle pastures and domestic crops created by modern, multi-purpose land use methods have put an unforseen burden on irrigation systems. The fences, necessary for separating cattle pastures from growing crops, greatly hinder the mobility of irrigation systems over a tract of multi purpose farmland. Large irrigation systems designed to efficiently irrigate large tracts of land, such as a pivotal irrigation systems, become impractical if they are obstructed by fences.\nOne approach to allow a fenced area to be irrigated, has been to place gaps in the fence which are wide enough for the cleated wheels of the irrigation system to pass. Unfortunately, the cleated wheels are so wide that the gaps in the fence are large enough to let livestock, including cattle, pass as well.\nAnother approach is to have crews move the fence in the path of the irrigation system and then replace it. This is labor intensive and expensive. It would therefore be advantageous to have fencing system that could articulate under the force of the moving irrigation system wheels, yet be resilient enough to retain livestock. Unfortunately, there are a number of problems associated with such a system. First, the fence must be able to yield to the force of the cleated wheels and articulate from its upright position, perpendicular to the ground, to a substantially flattened position that is parallel with the ground. Second, the fence system must maintain the strand integrity to keep from snapping wire strands as the fence articulates.\nArticulating fence systems have been suggested for varying uses, such as flood-fences built along river banks that are prone to regular flooding. U.S. Pat. No. 567,333 issued to Chandler, for example, discloses fence posts and fences that can pivot on a hinge from an upright position normal to the ground, to a flattened position parallel with the ground. Unfortunately, a number of problems arise when attempts are made to adapt these fence systems to use with mobile irrigation systems on farmland.\nMany hinged fence systems disclosed in the prior art are designed to remain parallel with the ground after being knocked down, and must be re-erected by hand to their upright positions. Such systems defeat the benefit of having a computer-controlled, self-propelled irrigation system, since they must be monitored and re-erected. Moreover, fences which are down for even a short time create a possibility that livestock will migrate from their designated pastures. In addition, they are labor intensive and thus expensive.\nThere are additional problems with fence systems disclosed in the prior art. Those that have some kind of resilient means, such as a spring mechanism, which allow the fence to spring back to an upright position after being knocked down by the irrigation system, are hinged at their base so as only to pivot in a single plane. If an irrigation system does not approach the fence in a direction that is perfectly aligned with the articulation direction of the fence system, then the irrigation system imparts side loads on the fence and can cause permanent damage to both systems. This problem is particularly acute for the widely used irrigation systems that rotate around a fixed center point. These systems commonly exert side load forces on the fence, causing either the posts or the wire strand to break or to be damaged.\nPresently, there is a need for an articulating fencing system that can allow passage of the irrigation system and accommodate the side loads exerted on the fence by the irrigation system. This need is especially great for fence systems which have delicate strands of electrified fencing wire which are easily snapped under the strain of an impinging irrigation system that moves along an arc that is not perpendicular to the fence line.\nIt has now been discovered that the problems encountered with prior art fencing systems can be overcome by the present invention. In the broad aspect of the present invention, a fencing system is provided that comprises at least two fence posts, wherein at least one of the fence posts is a resilient, isotropically articulating member, and the posts have strung there between at least one strand of flexible fencing material, such as a wire or cord. In accordance with the present invention, the system resiliently yields to forces exerted in substantially all directions, including directions other than those perpendicular to the fence line.\nIn accordance with the broad aspect of the present invention, the fencing system contains at least one isotropic fence post that resiliently yields to an object (e.g., a mobile irrigation system) approaching the fencing system from any direction, even a direction substantially parallel to the fence line. The fencing system yields to objects resistively. For example, the system will resiliently yield to a mobile irrigation system, but not to the force of a cow moving against the fence. In a preferred aspect, the resistance to a force pressing against the fence system is variably set, causing the fence posts to substantially yield only when a preset amount of force is applied. For example, the resistance can be set such that the force of a strong wind or a herd of cows will not cause the fence posts to articulate, but the force of an impinging irrigation system will.\nPreferably, the fence system of the present invention has at least three posts, including at least one resiliently articulating center post that is aligned between a pair of end posts, wherein a flexible fencing material is strung between adjoining posts. The at least one center post of the present invention is a resilient, isotropically articulating member that preferably includes a shaft adapted for affixing flexible fencing material; a resilient, isotropically flexible member having a top end for rigid communication with the shaft and a bottom end for rigid communication with a base element; and an anchor element, connected to the base element, for anchoring the base. The flexible fencing material preferably comprises strands of wire, either single strand or braided, including strands of electrically conducting metal for electric fences, and strands of barbed wire.\nThe end post elements of the system of the instant invention can be resiliently articulating, or they can be rigid, non-articulating supports. In a preferred aspect, the end posts are non-articulating. It will be realized that in accordance with the invention the end post elements may also form the apex of a corner that is formed from two linear segments of fence which are joined together at an angle (i.e., the end post is a corner post). At least one segment so joined is radially articulating.\nIn accordance with a preferred aspect of the present invention, the resilient, isotropically articulating posts have an isotropically flexible member for allowing a shaft to resiliently articulate between a position substantially perpendicular to the ground, and a position substantially parallel to the ground. In a preferred embodiment, the isotropically articulating member is a spring, and preferably a coil spring, that is able of radial, resilient articulation. In another embodiment, the isotropically flexible member is a segment of flexible hose, preferably made from plastic or rubber.\nPreferably, a fence material tensioning means, such as a coil spring, is used to attach a strand of fencing material to at least one post element of the fencing system. The tensioning means provides flexibility to assure the integrity of the fence material when the fence is under stress. This is to further assure that strands of fencing material in the fencing system will not break when the fence posts are articulated.\nIn one preferred embodiment, the tensioning means includes at least one spring element connected at one end to a strand of fencing material and on the other to a variable set tensioning means, such as a turnbuckle, which adjusts the amount of tension on the strand. In another embodiment, the tensioning means comprises two spring segments having a variable set tensioning means there between. In another embodiment, the tensioning means comprises a spring loaded pulley assembly rotatably attached to the end post assembly. A strand of fencing material is spooled on the pulley and held in tensioned engagement therewith."}
{"text": "Agents acting at central cholecystokinin (CCK) receptors may induce satiety (Schick, Yaksh and Go, Regulatory Peptides 14:277-291, 1986). They are also expected to act as analgesics (Hill, Hughes and Pittaway, Neuropharmacology 26:289-300, 1987), and as anticonvulsants (MacVicar, Kerrin and Davison, Brain Research 406:130-135, 1987).\nReduced levels of CCK-peptides have been found in the brains of schizophrenic patients compared with controls (Roberts, Ferrier, Lee, Crow, Johnstone, Owens, Bacarese-Hamilton, McGregor, O'Shaughnessey, Polak and Bloom. Brain Research 288:199-211, 1983). It has been proposed that changes in the activity of CCK neurones projecting to the nucleus accumbens may play a role in schizophrenic processes by influencing dopaminergic function (Totterdell and Smith, Neuroscience 19:181-192, 1986). This is consistent with numerous reports that CCK peptides modulate dopaminergic function in the basal ganglia and particularly the nucleus accumbens (Weiss, Tanzer, and Ettenberg, Pharmacology, Biochemistry and Behaviour 30:309-317, 1988; Schneider, Allpert and Iversen, Peptides 4:749-753, 1983). It may therefore be expected that agents modifying CCK receptor activity may have therapeutic value in conditions associated with disturbed function of central dopaminergic function such as schizophrenia and Parkinson's disease.\nCCK and gastrin peptides share a common carboxy terminal pentapeptide sequence and CCK peptides can bind to the gastrin receptor of the stomach mucosa and elicit acid secretion in many species including human (Konturek, Gastrointestinal Hormones, Ch. 23, pp 529-564, 1980, ed. G. B. J. Glass, Raven Press, NY). Antagonists of the CCK-B receptor would also be expected to be antagonists at the stomach gastrin receptor and this would also be of value for conditions involving excessive acid secretion.\nCCK and gastrin peptides have trophic effects on the pancreas and various tissues of the gastrointestinal tract (Johnson, ibid., pp 507-527), actions which are associated with increased DNA and RNA synthesis. Moreover, gastrin secreting cells are the Zollinger-Ellison syndrome (Stadil, ibid., pp 729-739), and some colorectal tumors may also be gastrin/CCK dependent (Singh, Walker, Townsend and Thompson, Cancer Research 46:1612, 1986; and Smith, J. P., Gastroenterology 95:1541, 1988). Antagonists of CCK/gastrin receptors could therefore be of therapeutic value as antitumor agents.\nThe CCK peptides are widely distributed in various organs of the body including the gastrointestinal tract, endocrine glands, and the nerves of the peripheral and central nervous systems. Various biologically active forms have been identified including a 33-amino acid hormone and various carboxyl-terminus fragments of this peptide (e.g., the octapeptide CCK26-33 and the tetrapeptide CCK30-33). (G. J. Dockray, Br. Med. Bull. 38(3):253-258, 1982).\nThe various CCK peptides are thought to be involved in the control of smooth muscle contractility, exocrine and endocrine gland secretion, sensory nerve transmission, and numerous brain functions. Administration of the native peptides cause gall bladder contraction, amylase secretion, excitation of central neurons, inhibition of feeding, anticonvulsive actions and other behavioral effects. (\"Cholecystokinin: Isolation, Structure and Functions,\" G. B. J. Glass, Ed., Raven Press, New York, 1980, pp 169-221; J. E. Morley, Life Sciences 27:355-368, 1980; \"Cholecystokinin in the Nervous System,\" J. de Belleroche and G. J. Dockray, Ed., Ellis Horwood, Chichester, England, 1984, pp 110-127.)\nThe high concentrations of CCK peptides in many brain areas also indicate major brain functions for these peptides (G. J. Dockray, Br. Med. Bull. 38(3):253-258, 1982). The most abundant form of brain CCK found is CCK26 33, although small quantities of CCK30-33 exist (Rehfeld and Gotterman, J. Neurochem. 32:1339-1341, 1979). The role of central nervous system CCK is not known with certainty, but it has been implicated in the control of feeding (Della-Fera and Baile, Science 206:471-473, 1979).\nCurrently available appetite suppressant drugs either act peripherally, by increasing energy expenditure (such as thyroxine), or in some other manner (such as the biguanides), or act by exerting a central effect on appetite or satiety.\nCentrally acting appetite suppressants either potentiate central catecholamine pathways and tend to be stimulants (for example, amphetamine), or influence serotonergic pathways (for example, fenfluramine). Other forms of drug therapy include bulking agents which act by filling the stomach, thereby inducing a \"feeling\" of satiety.\nCCK is known to be present in some cortical interneurones which also contain gamma-aminobutyric acid (GABA) (H. Demeulemeester et al, J Neuroscience 8, 988-1000, 1988). Agents that modify GABA action may have utility as anxiolytic or hypnotic agents (S. C. Harvey, The Pharmacological Basis of Therapeutics (7th ed.) 1985, pp 339-371, MacMillan). Thus, agents which modify CCK action may have parallel anxiolytic or hypnotic activities. The role of CCK in anxiety is disclosed in TIPS 11:271-273, 1990."}
{"text": "A quantum key distribution system is configured with a transmitter, a receiver, and an optical fiber link that connects the transmitter and the receiver. The transmitter transmits photons to the receiver via the optical fiber link (a quantum communication channel). After that, the transmitter and the receiver exchange control information with each other, and share cryptographic keys. This technology is implemented using the technology generally referred to as quantum key distribution (QKD).\nIn order to share cryptographic keys between a transmitter and a receiver using quantum key distribution, it is necessary to perform a key distillation operation in the transmitter as well as in the receiver. The key distillation operation includes a sifting operation, an error correction operation, and a privacy amplification operation. As a result of performing the key distillation operation, the transmitter and the receiver share cryptographic keys. The amount of generation per unit time of the shared cryptographic keys is called a secure key rate. Being able to use a number of cryptographic keys enables performing high-speed and safer cryptographic data communication. Hence, it can be said that, higher the secure key rate, the more enhanced is the performance of a quantum key distribution system.\nIn such quantum key distribution systems, there is a system in which an optical processing device, a high-speed signal processing unit, and a central processing unit (CPU) are arranged in series for the purpose of quantum key distribution and, depending on the details of communication to be done with another communication device, switching is done between whether the key distillation operation is to be assigned to the high-speed signal processing unit or to the CPU. In this way, the processing load and the communication load required in the key distillation operation is distributed to achieve a high speed during the operations.\nHowever, in such a quantum key distribution system, the optical processing device, the high-speed signal processing unit, and the CPU are connected in series; and the operations include sequential transfer of data. Hence, in case some of the operation modules (for example, the high-speed signal processing unit) stop operating, then the operations at the prior stage and the subsequent stage also get terminated. Moreover, the key distillation operation includes a plurality of different algorithms intended for the sifting operation, the error correction operation, and the privacy amplification operation. Hence, if the configuration includes only a single high-speed signal processing module, then it is not possible to have a high-speed signal processing unit with the most suitable configuration for each operation constituting the key distillation operation."}
{"text": "Milling tools are known for use in removing a section or window of existing casing from a well bore. In a known implementation, a whipstock-mill combination tool is run into the wellbore. This combination tool typically comprises a milling bit which is secured to a top of the whipstock. The milling bit is in some instances a starting mill. At a desired wellbore depth, typically set with a packer, the whipstock (which is hanging from the milling bit) is set at a correct orientation and secured in place. The milling bit is then freed from the whipstock and rotated with applied pressure. The concave face and taper of the whipstock forces the milling bit against the inside of the casing to cut an initial casing window (this initial milling operation may also remove the mechanism which attached the top of the whipstock to the milling bit). The starting mill is then removed from the wellbore and a window mill bit, attached to a more flexible drill string, is lowered into the wellbore. This window mill bit is then rotated to mill down from the initial casing window as the window mill bit rides further down the concave face and taper of the whipstock. This forms a casing window opening. The window mill bit is then removed. Additional mill bits may be run in the wellbore to define the size and shape of the opening in the casing. To the extent additional formation drilling is desired after opening the casing window, for example in connection with directional drilling operations, a drill bit would then be run into the wellbore and out through the formed casing window to engage the formation.\nIt will be recognized by those skilled in the art that the operation for milling a casing window is quite time consuming and expensive since it requires several different tools and multiple trips down the wellbore. If further formation drilling is required after opening the casing window, yet another trip is needed to run the formation drilling equipment back into the wellbore. There would be an advantage if more efficient milling and drilling methods and tools were available for use in connection with the milling of a casing window and the drilling of formations through that casing window.\nIn connection with known one trip mill-drill systems, which utilize a whipstock hung under the mill or mill drill for whipstock run in, orientation, and placement, the typical method for attaching the whipstock to the mill or mill drill utilizes one or more threaded sockets or a threaded nozzle port in the lower portion of the bit. An intervening component is then bolted to the threaded socket(s) and to the upper portion of the whipstock assembly. Operation of the system involves shearing off the bolted connection with the application of weight and rotational force. This weight and rotational force, through the concave face and taper of the whipstock, further urges the mill drill into the casing wall so as to commence window milling.\nFor efficiency and trip savings reasons, tool designers try to design the mill drill to be capable of both milling out the casing window and effectively drilling ahead into and through the rock formations beyond the casing. This design goal has not generally been satisfactorily achieved. One reason for this is that the diamond layers of PDC drill bit cutters which tool designers prefer for use in effectively removing formation material are not as effective for use in milling steel casing material. It is known that other materials, such as tungsten carbide, or cubic boron nitride (CBN), are better at cutting ferrous materials, like well casings, but are not as effective at cutting rock that is encountered for instance after the casing has been penetrated. To address this issue, it is known in the art to provide PDC drill bits with tungsten carbide, or cubic boron nitride (CBN), cutting features to assist in milling operations, while the PDC cutters of the bit function in connection with subsequent formation drilling. The performance of these combination mill/drill bits has not been entirely satisfactory.\nAn issue also arises with respect to connecting the whipstock to a PDC drill bit. Most PDC drill bits used in oilfield applications today are manufactured with bodies cast from a tungsten carbide matrix. It is extremely difficult and time consuming to attempt to retrofit a matrix PDC bit design with a threaded socket for whipstock attachment because the matrix material is too hard for standard machining. Alternatively, the threaded socket can be included in the original casting design. However, the design and integrity of the bit can be compromised by the inclusion of the threaded socket (either machined in, or cast in). These problems exist as well for bits whose bit bodies are made from material other than tungsten carbide. The incorporation of threaded sockets for whipstock attachment in steel bit bodies, for example, may also compromise optimum design or bit body integrity.\nThere accordingly exists a need in the art for an attachment mechanism for hanging a whipstock from a drag-type PDC drill bit. Preferably, a non-invasive way of attaching a whipstock to a standard drag-type PDC drill bit would be provided. This solution could also be applied to application specific mills or mill drills so as to eliminate the necessity for an invasive threaded socket to attach the whipstock.\nReference is made to prior art U.S. Pat. Nos. 7,178,609, 5,887,655, 3,652,138 and 5,069,297, the disclosures of which are hereby incorporated by reference in their entirety. Reference is also made to the Knight Fishing Services, X-1 Single Trip Whipstock; the Smith Services Trackmaster Plus Wellbore Departure System, the Weatherford Quickcut Casing Exit System, and the Western Well Tool Non-Rotating Protectors, prior art devices (incorporated by reference)."}
{"text": "Building panels provided with a mechanical locking system comprising a displaceable and resilient tongue cooperating with a tongue groove for vertical locking is known and disclosed in, e.g., WO2006/043893. The tongue is a separate part and is made of e.g. plastic and inserted in a displacement groove at an edge of a panel. The tongue is pushed into the displacement groove during a vertical assembling of the panels and springs back into the tongue groove of an adjacent panel when the panels have reached a locked position.\nAlso known is a locking system for panels comprising a tongue, which is displaceable along the edge of a panel, see, e.g., WO2009/116926, and cooperates with a tongue groove for vertical locking. The tongue is a separate part and is provided with several protrusions, which initially match recesses of the tongue groove. The panels may be assembled by a vertical movement and the tongue is displaced to a position in which the protrusions no longer match the recesses in order to obtain the vertical locking.\nAlthough the description relates to floor panel, the description of techniques and problems thereof is applicable also for other applications, such as panels for other purposes, for example wall panels, ceiling panels, furniture etc.\nA drawback with the known system is that a separate tongue must be produced and special inserting machines are required to position the tongue in the displacement groove with high precision.\nThe above description of various known aspects is the applicant's characterization of such, and is not an admission that any of the above description is considered as prior art."}
{"text": "1. Field of the Invention\nThe invention relates to quick release devices for constrained animals. In particular, the invention relates to a quick release mechanism incorporated in an animal constraining line such as a horse cross-tie or a dog leash wherein one end of the line is attached to the animal and the other end is constrained.\n2. Related Art\nIt is known to incorporate various quick release mechanisms into devices used to constrain or handle animals. Quick release mechanisms currently in use include biased latch-type mechanisms which can easily become corroded or contaminated by moisture, dirt, and debris, thereby making their operation unreliable.\nOne type of release which does not include a biasing member, pivoting member, or camming-action member is a xe2x80x9cthree-ringxe2x80x9d type release which is of a type used in the parachute art. However, to the best of Applicant\"\"s knowledge, an animal constraint or handling device effectively able to incorporate such a release has yet to be developed.\nSpecific restraints or handling devices which could benefit from a more reliable quick release mechanism include cross-ties and dog leashes.\nAs regards dog leashes, the handlers of guard dogs, attack dogs and the like (e.g., security personnel) generally constrain the dogs by holding onto a leash attached to a dog\"\"s harness, collar, etc. There may come a time when it is necessary to release the dog as rapidly as possible. In such an event, the need for the handler to manipulate a buckle or clip holding the leash to the collar or harness is time consuming, especially if the dog is straining at the leash and moving about. Also, it may be necessary for the handler to put down objects that he/she is carrying, such as a weapon or flashlight, or take his/her eyes from a subject, thereby placing the handler at a disadvantage. Moreover, when it later becomes necessary to reattach the dog to the leash, the handler may again be placed at a disadvantage if the reattachment process is difficult.\nQuick release dog leashes that can be remotely operated are disclosed, for example, in U.S. Pat. No. 5,125,365 and 4,149,492. However, these mechanisms tend to be of complicated construction, e.g., they may require springs, pivots, and/or cams, and are thereby expensive and prone to contamination and corrosion.\nAs regards cross-ties, such devices are used to restrain horses wherein one end of a cross-tie is attached to a harness or the like, and another end is held, e.g., by being attached to a fixed object. It may occur that it is desirable to release the horse as rapidly as possible if the horse has lost its balance or otherwise put itself in a precarious position.\nU.S. Pat. No. 962,890 to Byrd et al. discloses a harness for controlling animals. A release mechanism is provided for the xe2x80x9coperable ropesxe2x80x9d of the restraint. The release is located along the lower side area of the animal, which is relatively difficult to access when the animal is standing. The release is activated by pressing a thumb piece, which requires the handler to position himself close to the animal in order to activate the release mechanism, thereby increasing the danger the handler will be kicked or otherwise harmed by the animal. Also, the release utilizes spring-biased mechanical parts which are susceptible to contamination, corrosion, and subsequent failure as noted above.\nU.S. Pat. No. 5,197,410 to Wilson et al. discloses a releasable tether having a hook-shaped end and a spring-biased bolt closing member. The release mechanism includes a cable housed in a plastic tube. One end of the cable is connected to the bolt, the other end is connected to a trigger member. This release mechanism is rather complex, and includes moveable spring-biased parts which are prone to contamination, corrosion, and subsequent failure.\nU.S. Pat. No. 5,103,771 to Lee discloses a remotely controlled quick release device. The release can be actuated by electromechanical means or manually by purely mechanical means. The release disclosed in Lee is overly complicated, and includes spring-biased relatively moveable parts that are subject to corrosion and contamination as noted above.\nIt is an object of the present invention to provide improved quick release animal constraints and handling devices that are simple, reliable, and inexpensive.\nIt is another object to provide a quick release dog leash which provides a way of quickly re-attaching the leash to a previously released dog.\nThese and other objects are attained by providing an animal constraint or handling device with a quick release mechanism that is essentially free of any biasing member, pivoting member, or camming-action member.\nA cross-tie constructed according to the present invention includes first and second end cross-tie portions connected by a quick release mechanism that is essentially free of any biasing member, pivoting member, or camming-action member.\nA plurality of cross-ties constructed according to the present invention each include a quick release mechanism and an associated trigger mechanism which is slidably connected to a respective cross-tie. The plurality of trigger mechanism are connected to each other and form a singular grasping portion, thereby permitting activation of the plurality of quick release mechanisms of the plurality of cross-ties via the singular grasping portion.\nThe present invention relates to a lead for constraining an animal, the lead comprising a line having a first end portion adapted to be connected to an animal, and a second end portion adapted to be constrained. The lead further includes a coupler connectable to the animal, and a quick release mechanism for connecting the coupler to the first end portion of the line. The quick release mechanism is essentially free of any biasing member, pivoting member, or camming-action member. The line further includes a trigger member operably slidable relative to, and extending away from, the quick release mechanism. The trigger member is manually graspable at a location spaced from the quick release mechanism, thereby enabling remote actuation of the quick release mechanism.\nPreferably, the lead constitutes a dog leash, and the coupler constitutes a first coupler. The lead further comprises a second coupler attached to the line for re-leashing a previously released dog.\nThat second coupler is preferably held against relative movement to the line by the quick release mechanism."}
{"text": "The invention relates to a regional navigation satellite supplementary system.\nA navigation system must not only allow positioning of the user with the required precision, continuity, and availability, but also provide the other information necessary for navigation. This other information is typically provided on paper via a map, on a storage media, via traffic radio messages, which are also distributed in digital form.\nThe street grid maps used in the navigation systems are usually fixed and typically require user interaction to obtain regular updates. Using known techniques, information about current or future traffic obstructions is available only on the local or regional level (if at all), and interregional information (which is also significant for planning longer trips) is typically not easily available. Accordingly, conventional systems do not necessarily provide selective detour recommendation that account for new traffic obstructions that could arise due to bypass recommendations.\nExemplary embodiments of the present invention provide devices and methods which allow traffic obstructions to be bypassed. This allows the street infrastructure to be used significantly more effectively because traffic obstructions may be bypassed. Moreover, the present invention significantly reduces the effort that a user must expend to have current planning documents for his trips.\nIn accordance with exemplary embodiments of the present invention, the navigation satellites of the navigation satellite supplementary system which are located over an area in which a service is to be provided are located close to the zenith and transmit a relatively strong navigation signal into the service area. For example, these satellites may fly on suitable geosynchronous orbits, on geosynchronous orbits in or near the equatorial level, or also on so-called Molniya orbits.\nThese satellites are then implemented as controllable using a (vertically) oriented antenna, to achieve the highest possible signal strength in the service area with the least possible power consumption, the signal strength being dimensioned in such a way that the signals of other satellite navigation systems may still be received. This is achieved, inter alia, by a suitable selection of the coding of the signals.\nThe signals transmit navigation-relevant information at a data rate which is significantly above the data rate of the commercial service of Galileo. This information contains, inter alia, current and expected traffic obstructions in the service area, recommendations for how specific user groups having various travel directions and travel destinations (which may be regionally combined) are to drive around these traffic obstructions, changes of the street grid in the service area, and also comprise expected traffic obstructions due to environmental influences, chronologically recent changes being repeated relatively frequently and chronologically earlier changes being repeated relatively rarely (which may be no repetition at all). The invention also discloses arbitrary combinations of the above-mentioned navigation-relevant information.\nExemplary embodiments of the present invention include a method for transmitting a navigation signal that involves generating navigation signals, wherein the signal strength of the navigation signals is dimensioned in such a way that it is still possible to receive signals of other satellite navigation systems, and wherein the signals transmit navigation-relevant information at a data rate which is significantly above the data rate of the Galileo satellite system; and transmitting the generated navigation signals.\nThe features specified in connection with the embodiments of the device of the invention are also disclosed for the method.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of one or more preferred embodiments when considered in conjunction with the accompanying drawings."}
{"text": "1. Technical Field\nThe present disclosure relates to an adjustable bed system, and more particularly, to an adjustable bed system with a bed frame that is adjustable in height.\n2. Background of Related Art\nAdjustable beds are often used in both home care, and in more formalized medical settings, e.g., hospital rooms. Adjustable beds generally include a pair of end boards, i.e., a headboard and a footboard, a bed frame that extends between the end boards to support a mattress, and a mechanism that allows the height of the bed frame to be adjusted between the end boards so that the bed frame, and thus the mattress and patient, can be raised and lowered.\nVarious height adjustment mechanisms are known in the art, and typically include a pair of transition boxes, or gearboxes, that are positioned on the end boards, i.e., one transition box on the footboard, and another transition box on the headboard. The transition boxes include internal gearing mechanisms, and are connected to drive screws extending vertically through the end boards such that upon actuation of the transition boxes, the drive screws rotate to either raise or lower the bed frame dependent upon the direction of rotation. One example of such an arrangement is described in U.S. Pat. No. 5,134,731 (hereinafter “the '731 patent”).\nAdjustable bed systems can be either manually operated, or automatic. Manual systems utilize transition boxes that are operated via a hand crank, for example, whereas automated systems regulate operation of the transition boxes via an electric motor. In both manual and automated systems known in the art, the transition boxes are arranged on the end boards so that they face each other when the system is assembled. A drive shaft extends between, and connects, the transition boxes so that the actuation of one transition box causes corresponding actuation of the other. More specifically, since the drive shaft is connected to both the transition boxes, actuating one of the transition boxes causes rotation of the drive shaft, which thereby transmits a rotational force to the other transition box to the cause simultaneous actuation.\nIn adjustable bed systems such as that described in the '731 patent, the end boards are different, in that the transition boxes included on the headboard and the footboard are configured for rotation in opposite directions during use. However, such systems have led to inefficiencies during delivery and assembly. For example, on the occasion that two headboards or two footboards are inadvertently delivered, as opposed to one headboard and one footboard, the system would not function properly upon assembly, if at all. In order to remedy the predicament, the bed system would have to be disassembled, and the appropriate parts, i.e., either the missing headboard or footboard, would have to be re-delivered, resulting in not only increased operational costs, but customer dissatisfaction as well.\nSystems such as those described in U.S. Pat. Nos. 6,983,495, 6,997,082, 7,302,716, and U.S. Pat. No. 7,441,289 have attempted to prevent such delivery and assembly issues via the development of identical headboards and footboards. Utilizing identical headboards and footboards reduces manufacturing costs, while also eliminating the chance for delivery of an improper end board. These systems, however, are incompatible with systems such as those described in the '731 patent.\nAccordingly, the present disclosure is directed to an improved adjustable bed system, and in particular, to an improved bed frame, that is universal in the sense that it can be used with different end boards, such as those described in the '731 patent, as well as with identical end boards, such as those described in U.S. Pat. Nos. 6,983,495, 6,997,082, 7,302,716, and U.S. Pat. No. 7,441,289."}
{"text": "The integration of high value capacitors in integrated circuits (ICs) is limited by the fact that conventional high value capacitors take up large areas of the IC chip, thus reducing device packing density and layout efficiency. Many applications require a large number of capacitors. Often the capacitors must be incorporated as discrete off-chip components, substantially increasing the bulk of the peripheral circuitry. In view of the increasing demand for compact lightweight electronic equipment, it is desirable that the number of discrete components be reduced.\nThe minimum dimensions of IC capacitors are determined primarily by the relatively low dielectric constant (.di-elect cons.&lt;10) of conventional capacitor dielectrics, e.g. SiO.sub.2 and Si.sub.3 N.sub.4. Thus, as device dimensions decrease, there is increasing interest in other dielectrics having higher dielectric constants than conventional dielectric materials.\nFerroelectric (FE) and high-epsilon (HE) dielectrics (.di-elect cons.=20 or greater) have found application in random access memory (RAM) cells since they provide for the formation of NVRAMs and DRAMs. Moreover, when used as NVRAM, ferroelectric dielectrics advantageously provide low voltage programmability, fast access times and low power consumption.\nFE materials pose several integration problems. In particular, most ferroelectric materials require high temperature post-deposition oxygen anneals (600.degree. C. or above) to achieve properties desirable for storage media. However, such high anneal temperatures can be incompatible with the CMOS devices already fabricated on the wafer. Furthermore, any subsequent forming gas or hydrogen anneals (highly desirable for CMOS devices) degrade the ferroelectric material, thus requiring additional high temperature oxygen anneals late in the processing sequence which in turn are detrimental to the CMOS circuitry.\nIn view of the above drawbacks with prior art methods of integrating ferroelectric capacitors with CMOS structures, there is a continued need for developing a new and improved method which is capable of providing an integrated FE capacitor/CMOS structure without subjecting the CMOS structure to high temperature steps that are typically required in the prior art for ferroelectric processing."}
{"text": "The invention relates to the sequential storage circuits commonly referred to as xe2x80x9cflip-flopsxe2x80x9d. More particularly, the invention relates to a high-speed flip-flop that operates at very low voltage levels and can offer set and/or reset capability.\nFlip-flops are sequential circuits storing either a xe2x80x9chighxe2x80x9d value (power high, or logic one) or a xe2x80x9clowxe2x80x9d value (power low, or logic zero). A flip-flop has a next value that depends on the values of one or more input signals. Conventionally, a flip-flop has data, clock, set, and/or reset input signals.\nA D (data) input signal is typically clocked into the flip-flop on receipt of a given clock edge, and appears at the flip-flop output on the opposite clock edge. S (set) and R (reset) input signals are generally unclocked, meaning that when the set or reset signal becomes active (e.g., goes high), the stored value changes immediately, without waiting for the arrival of a clock edge. An active set signal forces the stored value (conventionally designated Q) high, no matter what value was previously stored. An active reset signal forces the stored value Q low, no matter what value was previously stored. In set/reset flip-flops (i.e., flip-flops having both set and reset input signals) the set and reset signals are typically restricted such that at most one of them can be active at any given time.\nFlip-flops are often designed using two latches separated by passgates. FIG. 1 shows such a flip-flop, comprising a first latch including cross-coupled inverters 101 and 102, a second latch including cross-coupled inverters 103 and 104, and passgates 106, 107. The clock signal C is inverted by inverter 105 to provide inverted clock signal CB. On the falling edge of the clock signal C (or whenever signal C is low), signal CB goes high, and the data signal D passes through passgate 106 to node A, and hence into the first latch. (In the present specification, the same reference characters are used to refer to terminals, signal lines, and their corresponding signals.) The data value is inverted by inverter 101, and appears at node B. On the next rising edge of clock signal C, the inverted data from node B passes through passgate 107 to node QB (the inverted output signal), and hence into the second latch. The data is again inverted by inverter 103, and appears on the output node Q.\nFor the flip-flop to function properly, it must be possible to write a high value to each of nodes A and QB, and ensure that the new value overcomes a low value previously stored on the node. The new high value must pass through either passgate 106 or passgate 107, and overcome the zero value being provided by inverter 102 or inverter 104, respectively. When passing through an N-channel transistor, a high voltage value is reduced by the threshold voltage of the transistor. Therefore, if an N-channel transistor is used to implement the passgate, as shown in FIG. 1, inverters 102 and 104 are necessarily designed to be weak, and correspondingly slow. (Writing low values is not an issue, because there is no voltage degradation of a low signal through an N-channel transistor.)\nAt very low voltages, the high voltage value (VDD) approaches the threshold voltage of an N-channel transistor, which is generally about 0.7 volts. It is not practical to reduce this threshold voltage, because transistors with a lower threshold voltage would become unacceptably sensitive to noise. Therefore, at very low voltages, the high value passed through the passgate is not sufficient to overcome even a weak inverter (e.g., inverter 102 in FIG. 1), and therefore may not be high enough to trip inverter 101. The same limitation applies to passgate 107 and inverters 104, 103. One known solution is to implement the passgates using CMOS passgates (i.e., paired N-channel and P-channel transistors), as shown in FIG. 2.\nFIG. 2A shows a well-known flip-flop similar to that of FIG. 1. The N-channel transistors forming passgates 106, 107 in FIG. 1 are replaced by CMOS passgates 206, 207, respectively. The P-channel gate terminals are driven by the inverse of the signals used to drive the N-channel gate terminals. A high value on data input D or node B is not reduced as it passes through the P-channel transistor. Therefore, the substitution of the CMOS passgates for the simple N-channel transistors ensures that new high values will be successfully latched, even at low voltages. (Of course, the high voltage level must be higher than the threshold voltage of the N-channel transistor, or no circuit containing N-channel transistors will function properly.)\nA disadvantage of using CMOS passgates in flip-flops is that P-channel transistors are relatively slow compared to similarly-sized N-channel transistors. Because flip-flops are ubiquitous in nearly all integrated circuits, minimum sized gates are usually used to implement flip-flops. Passing a high logic level through a minimum sized P-channel transistor introduces a delay that can become significant.\nTherefore, the flip-flop of FIG. 2A also includes another modification, which improves the operating speed of the flip-flop. CMOS passgates 208 and 209 are inserted into the latch feedback loops after inverters 202 and 204, respectively. When clock signal C goes low and inverted clock signal CB goes high, CMOS passgate 206 is turned on. In the flip-flop of FIG. 1, an incoming high data signal must then overcome a low value on node A driven by inverter 102. In the flip-flop of FIG. 2A, however, passgate 208 is turned off (because clock signal C is low), and there is no driving inverter to be overcome by the new high value. Similarly, when clock signal C goes high and inverted clock signal CB goes low, passgate 209 is turned off, and a high value on node B passes onto node QB unopposed by inverter 204. Thus, the latches in FIG. 2A change state more rapidly than the corresponding latches in FIG. 1.\nFIG. 2B shows another prior art flip-flop similar to that of FIG. 2A. In the flip-flop of FIG. 2A, passgates (208, 209) are inserted between the outputs terminals of the feedback inverters (202, 204) in the latches and the nodes (A, QB) by which data is passed to the latches. These passgates provide a means for electrically isolating the feedback inverters from the-nodes when new data is provided to the latches. However, writing a low value to a latch is generally not a problem, even at low voltages. The problems are encountered only when writing a high value to the latch while overcoming a low value provided by the feedback inverter. Therefore, when writing any new value to the latch, it is only necessary to prevent the feedback inverter from driving a low value.\nIn the flip-flop of FIG. 2B, P-channel transistor 211 and N-channel transistor 213 form an inverter driven by node B (corresponding to inverter 202 in FIG. 2A). Inserted into the pull-down path of the inverter is an N-channel transistor 212 gated by the clock signal C. Thus, when clock signal C goes low, the new data (e.g., a high value) is written from data input terminal D to node A. Because clock signal C is low, no low signal is provided through N-channel transistor 213. Therefore, the new high value on node A is unopposed.\nSimilarly, P-channel transistor 221 and N-channel transistor 223 form an inverter driven by node Q (corresponding to inverter 204 in FIG. 2A). Inserted into the pull-down path of the inverter is an N-channel transistor 222 gated by inverted clock signal CB. Thus, when clock signal C goes high, the new data (e.g., a high value) is written from node B to node QB. Because inverted clock signal CB is low, no low signal is provided through N-channel transistor 223. Therefore, the new high value on node QB is unopposed.\nIt is possible to include feedback passgates 208 and 209 (FIG. 2A) or N-channel transistors 212 and 222 (FIG. 2B) while using N-channel transistors 106 and 107 (FIG. 1) instead of CMOS passgates 206 and 207 (FIGS. 2A and 2B) to feed the latches. Eliminating the feedback paths while writing to nodes A and QB gives improved performance compared to the flip-flop of FIG. 1. However, as described above, at very low voltages an incoming high value can be reduced below the threshold voltage of the N-channel device in the latch inverter. For example, if the power high voltage is 1.4 volts and the threshold voltage of an N-channel device is 0.7 volts, the high voltage after passing through N-channel transistor 106 or 107 is only 0.7 volts, which might or might not be sufficient to trip the latch. When processing variations are taken into account, the flip-flop becomes unreliable even at voltages higher than 1.4 volts. Therefore, the use of N-channel transistors to implement passgates 106 and 107 is not desirable for known flip-flop circuits expected to operate at low voltages.\nWhile CMOS passgates 206 and 207 are necessary to ensure functionality at low voltages, the P-channel transistors included in the passgates still introduce undesirable delay into the data path through the flip-flop when writing high values over low values previously stored in nodes A or QB.\nThe issue of speed becomes more problematical when reset and/or set capability is added to the flip-flop. FIG. 3 shows a well-known flip-flop similar to that of FIG. 2A, but having reset capability. The reset signal R is inverted through inverter 310 to provide inverted reset signal RB. Inverter 201 is replaced with NAND gate 301, and the second input terminal of NAND gate 301 is driven by inverted reset signal RB. Inverter 203 is replaced with NOR gate 303, and the second input terminal of NOR gate 303 is driven by reset signal R.\nThus, when signal R goes high, signal RB goes low and node B is forced to a high value, while the output signal Q is forced to a low value. Both values are latched through the respective feedback loops.\nThe reset flip-flop of FIG. 3 is slower than a corresponding flip-flop without reset capability, because of the extra transistors included in NAND gate 301 and NOR gate 303. Assume that reset signal R is inactive (deasserted, or low) and signal RB is high. Whenever a high value is written from data input D to node A, node B goes low. However, node B is now pulled low through two N-channel transistors in series. (See FIG. 3A, which shows a standard CMOS NAND-gate that includes two N-channel transistors 313, 314 coupled in series and two P-channel transistors 311, 312 coupled in parallel.) Therefore, NAND gate 301 is slower than inverter 201.\nSimilarly, whenever a low value is written from node B to node QB, output node Q goes high (assuming the reset signal R is inactive). However, output node Q is now pulled high through two P-channel transistors in series. (See FIG. 3B, which shows a standard CMOS NOR-gate that includes two P-channel transistors 315, 316 coupled in series and two N-channel transistors 317, 318 coupled in parallel.) Therefore, NOR gate 303 is slower than inverter 203.\nFIG. 4 shows a well-known flip-flop similar to that of FIG. 2A, but having set capability. The set signal S is inverted through inverter 460 to provide inverted set signal SB. Inverter 201 is replaced with NOR gate 401, and the second input terminal of NOR gate 401 is driven by set signal S. Inverter 203 is replaced with NAND gate 403, and the second input terminal of NAND gate 403 is driven by inverted set signal SB.\nThus, when signal S goes high, node B is forced to a low value, signal SB goes low, and the output signal Q is forced to a high value. Both values are latched through the respective feedback loops.\nSimilar to the reset flip-flop of FIG. 3, the set flip-flop of FIG. 4 is slower than a corresponding flip-flop without set capability, because of the extra transistors included in NOR gate 401 and NAND gate 403.\nIt is part of the nature of transistors that, all else being equal, a lower voltage means a slower operating speed. At one time, virtually all integrated circuits operated at 5 volts. However, voltage levels of integrated circuits have been rapidly decreasing, first to 3.3 volts, then 2.5 volts, and so on. Voltage levels of 1.8 volts or less are now supported by many integrated circuits.\nAt the same time, the operating frequencies expected from the circuits have been increasing at a rapid rate. Where once 10 megahertz was considered an acceptable clock rate, clock rates of up to 300 megahertz are now commonly specified. Integrated circuit designers are being forced to redesign many common circuits to obtain additional performance at lower voltage levels.\nTherefore, it is desirable to provide high-speed flip-flops, particularly set and/or reset flip-flops, that function reliably at reduced voltage levels.\nThe invention provides flip-flops that are both operable at high speed and reliable at low voltage levels.\nAccording to a first aspect of the invention, a first flip-flop includes first and second latches. When a clock signal is in a first state, a data input signal is passed to a first node of the first latch through a first passgate, and the inverted data input signal is passed to a second node of the first latch through a second passgate. When the clock signal is in a second state, a signal on the second node of the first latch is passed to a first node of the second latch through a third passgate, and the inverse of that signal is passed to a second node of the second latch through a fourth passgate.\nIn one embodiment, the first, second, third, and fourth passgates are N-channel transistors. CMOS passgates are used in one embodiment. However, they are not necessary in most cases, because whenever a high value is-passed to one node of a latch, a low value is passed to the other node of the latch. Therefore, the latch can safely ignore all high input values. Hence, this feature permits the flip-flops of the invention to function at very low voltages. Because writing a high value is normally much slower than writing a low value, the flip-flops of the invention also function at very high clock rates, even at low voltages.\nIn some embodiments, means are provided for disabling the feedback loops of each latch while data is passed to the latch. The first latch comprises first and second logic gates, with the first logic gate driving the first node and the second logic gate, and the second logic gate driving the second node and the first logic gate. The second latch also comprises first and second logic gates, with the first logic gate driving the first node and the second logic gate, and the second logic gate driving the second node and the first logic gate.\nIn one such embodiment, a feedback passgate is coupled between an output terminal of the first logic gate and the first node of each latch. The N-channel gate terminal of the feedback passgate in the first latch is coupled to the gate terminal of the first passgate. The P-channel gate terminal (where provided) is coupled to receive the inverse signal from the N-channel gate terminal. The N-channel gate terminal of the feedback passgate in the second latch is coupled to the gate terminal of the second passgate. The P-channel gate terminal (where provided) is coupled to receive the inverse signal from the N-channel gate terminal.\nIn another such embodiment, an N-channel transistor is inserted into the pull-down path of the first logic gate of each latch. Thus, when the N-channel transistor is xe2x80x9coffxe2x80x9d, the logic gate does not pull the first node of the latch low, and a high value is easily written to the latch.\nA second aspect of the invention enables the use of cross-coupled inverters, rather than NAND gates and NOR gates, to implement set, reset, and set/reset flip-flops. Because inverters are faster than similarly-sized NAND and NOR gates, this features also increases the speed at which the flip-flop can operate.\nIn a flip-flop according to this aspect of the invention, a first latch comprises first and second inverters, with the first inverter driving the first node and the second inverter, and the second inverter driving the second node and the first inverter. The second latch also comprises first and second inverters, with the first inverter driving the first node and the second inverter, and the second inverter driving the second node and the first inverter. For each latch, at least one pull-up and at least one pull-down are coupled directly to the nodes of the latch. The pull-ups and pull-downs are controlled by set and/or reset signals applied to the flip-flop.\nIn a reset flip-flop according to one embodiment of the invention, a first pull-down is coupled to the first node of the first latch, with the gate terminal of the first pull-down being driven by a reset signal. A first pull-up is coupled to the first node of the second latch, with the gate terminal of the first pull-up being driven by an inverted reset signal. Optionally, a second pull-up is coupled to the second node of the first latch, with the gate terminal of the second pull-up being driven by the inverted reset signal, and a second pull-down is coupled to the second node of the second latch, with the gate terminal of the second pull-down being driven by the reset signal.\nIn a set flip-flop according to another embodiment of the invention, a third pull-up is coupled to the first node of the first latch, with the gate terminal of the third pull-up being driven by an inverted set signal. A third pull-down is coupled to the first node of the second latch, with the gate terminal of the third pull-down being driven by the set signal. Optionally, a fourth pull-down is coupled to the second node of the first latch, with the gate terminal of the fourth pull-down being driven by the set signal, and a fourth pull-up is coupled to the second node of the second latch, with the gate terminal of the fourth pull-up being driven by the inverted set signal.\nIn a set/reset flip-flop according to yet another embodiment of the invention, all of the above-mentioned first and second pull-ups and pull-downs are included. In another embodiment, all of the first, second, third, and fourth pullups and pull-downs are included. In another embodiment, only the third and fourth pull-ups and pull-downs are included. For all of the set/reset flip-flops of the invention, the set and reset signals are preferably restricted such that at most one of them can be active at any given time.\nAccording to a third aspect of the invention, both the first and second aspects are applied. In other words, when an input signal is applied to one terminal of each latch, an inverted input signal is applied to the other terminal of the same latch. In addition, set and/or reset capability is provided by coupling at least one pull-up and at least one pull-down directly to the nodes of the latch. Thus, set and/or reset flip-flops are provided having the advantages of both high speed and the ability to function properly at very low power levels."}
{"text": "1. Technical Field of the Invention\nThe invention is in the field of kinetic anti-projectile interceptor vehicles.\n2. Description of the Related Art\nInterceptors have been proposed to intercept and disable or destroy space-based or space-entering projectiles, for example ballistic projectiles such as intercontinental ballistic missiles. Such projectiles travel at very high rates of speed and have short travel times, making interception of them a difficult problem, one in which there is room for further improvements."}
{"text": "Non-aqueous electrolyte batteries that use a non-aqueous electrolyte such as lithium-ion secondary batteries provide a high energy density at a high voltage and have a reduced size and a reduced weight, and thus are widely used primarily as power sources for information terminals such as personal computers and cellular phones. Common examples of the non-aqueous electrolyte used in the non-aqueous electrolyte batteries include a solution obtained by dissolving a supporting salt such as LiPF6 in an aprotic organic solvent such as an ester compound and an ether compound. Because the aprotic organic solvent is flammable, however, the battery may be ignited when the liquid leaks from the battery or the battery generates an abnormal amount of heat, which is disadvantageous in terms of safety.\nThe possibility of expanded use of the non-aqueous electrolyte batteries as power sources for large devices such as power sources for power storage and power sources for electric vehicles has recently been studied. Therefore, it has been an important issue to enhance the safety of the non-aqueous electrolyte batteries with a view to increasing the size of the batteries. For example, Japanese Unexamined Patent Application Publication No. 2000-173619 (Patent Document 1) discloses a technology for suppressing ignition or rupture of the non-aqueous electrolyte batteries by covering the surface of a negative electrode with a phosphazene monomer which is a flame retardant material, as a technology for enhancing the safety of the batteries."}
{"text": "1. Field of the Invention\nThe present invention is directed to isolated switching voltage converters and more particularly to primary side sensing of the output voltage on the secondary side of a switching voltage converter.\n2. Description of the Related Art\nSwitching voltage converters (also referred to herein as switching voltage regulators) are used to provide, e.g., regulated DC output voltage from an unregulated AC input. Typical consumer products involving such switching regulators include cell phone chargers, laptop or printer power supplies (so-called “bricks”), and embedded PC power supplies.\nFIG. 1 illustrates a switching regulator 100 having a flyback topology commonly used in power supplies with output power less than 200 W. The illustrated topology is suitable for discontinuous current mode (DCM) operation in which the transformer core is fully demagnetized in each cycle. Another mode of operation is continuous current mode (CCM).\nThe switching regulator 100 includes high voltage isolation between the DC output (secondary side) and the AC mains input (primary side). The isolation is required for safety/regulatory reasons and may also be required for functional reasons. The power is transferred from the primary side to the secondary side using a transformer 102 meeting isolation requirements and primary side high voltage switch 104.\nThe switching regulator shown in FIG. 1 operates as follows for DCM operation. When the power MOSFET switch 104 turns ON (TON phase) according to the gate control signal supplied by the controller integrated circuit 106, the current through primary winding 105 ramps up with a slope of Vin/Lp and the energy stored in the transformer core at the end of the TON cycle is proportional to 0.5×Lp×Ippeak2, where Lp is the transformer primary winding inductance and Ippeak is the primary winding peak current. The output current Is is zero during the TON phase and the voltage Vs is negative referenced to the secondary side ground, Vs=−Ns*Vin, where Ns is the transformer secondary/primary turn ratio.\nWhen switch 104 turns OFF (TOFF phase), the primary inductor current Ip becomes zero and secondary current Is ramps down from the value Ispeak=Ippeak/Ns to zero, with a slope of approximately ˜(Vout+Vdout)/Ls. When Is>0 and Q1 is OFF, the output voltage is reflected according to the transformer turn ratio back to the primary side and is usually called primary flyback voltage: Vfly_p˜(Vout+Vdout)/Ns, where Vout is the output voltage, Vdout is a voltage drop across the output diode 107 and Ns is defined as above. Once the secondary winding current Is reaches zero, both transformer windings become open (provided that loading on auxiliary winding can be neglected) and flyback voltage converts to damped ringing waveform fueled by residual energy in the Lp and Cd resonant circuit, where Cd is the total equivalent capacitance at the drain of Q1.\nAccurate regulation of the output voltage requires feedback proportional to output voltage. The feedback controls the duty cycle of switch 104 in order to keep the output voltage constant over changing load and input voltage. The feedback path needs to cross the isolation barrier between the primary and secondary. A common feedback solution uses an opto-coupler 108 as shown in FIG. 1.\nThe use of an opto-coupler for feedback has costs such as additional components as well as PCB area space. An opto-coupler feedback approach typically requires additional passive and active components like resistors, capacitors and a shunt regulator. The shunt regulator 109 can be alternatively replaced by Zener diode when lower quality output voltage regulation accuracy can be tolerated.\nThere have been suggestions to replace the opto-coupler and regulate the output voltage by sensing the flyback voltage at either the primary winding or at an auxiliary winding (sometimes referred to as a bias winding). Such approaches are usually denoted as regulation using magnetic flux sensing or primary side sensing. However prior art primary side sensing approaches have had problems with accurately determining the flyback voltage. Accordingly, improved primary side sensing is desirable."}
{"text": "There is increasing demand for portable computing devices, such as a smartphone, adapted to perform a wide array of functions. The technological sophistication of these devices has increased to meet this demand, such that many portable computing devices are effectively mobile computers completed with an integral display. Because of their mobility and sophistication, said devices are particularly susceptible to damage. A number of cases have been developed to protect such devices from fall damage, water exposure, or like hazards.\nMany aspects of modern life have become dependent upon the use of portable computing or smartphone devices. For example, in addition to voice communication, many people have come to rely upon their smartphone for basic functions, such as e-mail and navigation. Numerous software applications have also been developed to monitor or diagnose certain health conditions using the basic functions of said devices, such as a camera and a microphone. Although useful, the specialized capabilities of many such applications are inherently limited by the capabilities of the smartphone itself.\nTherefore, continuing need exists for an electronic device holder that is adapted to protect an electronic device. It would be further desirable to provide a holder that can enhance the capabilities of said device."}
{"text": "Field\nThe present disclosure relates generally to a method and apparatus for detecting license plate information from an image of a license plate and more specifically, detecting license plate information from an optical image, captured by a mobile apparatus, that includes a license plate image and several other object images.\nBackground\nIn recent years, collecting still images of license plates has become a common tool used by authorities to catch the drivers of vehicles that may engage in improper or unlawful activity. For example, law enforcement authorities have set up stationary traffic cameras to photograph the license plates of vehicles that may be traveling above a posted speed limit at a specific portion of a road or vehicles that drive through red lights. Toll booth operators also commonly use such stationary cameras to photograph vehicles that may pass through a toll booth without paying the required toll. However, all of these scenarios have a common thread. The camera must be manually installed and configured such that it will always photograph the vehicle's license plate at a specific angle and when the vehicle is in a specific location. Any unexpected modifications, such as a shift in angle or location of the camera would render the camera incapable of properly collecting license plate images.\nAdditionally, camera equipped mobile apparatuses (e.g., smartphones) have become increasingly prevalent in today's society. Mobile apparatuses are frequently used to capture optical images and for many users serve as a replacement for a simple digital camera because the camera equipped mobile apparatus provides an image that is often as good as those produced by simple digital cameras and can easily be transmitted (shared) over a network.\nThe positioning constraints put on the traffic cameras make it difficult to take images of license plates from different angles and distances and still achieve an accurate reading. Therefore, it would be difficult to scale the same license plate image capture process performed by law enforcement authorities to mobile apparatuses. In other words, it is difficult to derive license plate information from an image of a license plate taken from a mobile image capture apparatus at a variety of angles, distances, ambient conditions, mobile apparatus motion, and when other object images are also in the image, which hinders a user's ability to easily gather valuable information about specific vehicles when engaging in a number of different vehicle related activities such as buying and selling vehicles, insuring vehicles, and obtaining financing for vehicles."}
{"text": "1. Field of the Invention\nThe present invention relates to a plate-like filter element, a filter unit formed by stacking up a plurality of the filter elements or a filter device installed therein, the filter unit being used to mainly trap or remove particulates contained in an exhaust gas discharged from a diesel engine which is mainly used for various vehicles, vessels or industrial appliances.\nAn exhaust gas discharged from a diesel engine contains particulates including carbon particles as a main component at a fairly high content, which causes air pollution.\n2. Description of the Related Art\nVarious devices have been proposed to trap or remove the particulates contained in the exhaust gas by using a filter means.\nFor instance, Japanese Unexamined Patent Publication No. 124417/1981 discloses a filter unit 20 made of ceramics as shown in FIG. 9. The filter unit 20 shows an outer configuration of a rectangular prism as a whole and comprises a plurality of (eight in FIG. 9) rectangular plate-like bodies 21, 22 placed in parallel to each other, vertical end ribs 23, 25 and a number of intermediate ribs 24, 26.\nAll the plate-like bodies 21, 22, ribs 23, 25 and intermediate ribs 24, 26 or at least the plate-like bodies 22 are made of gas-permeable porous thin ceramic wall having a filtering function. The plate-like bodies 21 constitute the upper and lower walls of the filter unit 20, and the plate-like bodies 22 constitute the intermediate walls. Between adjacent plate-like bodies 21, 22 and between adjacent plate-like bodies 22, 22, the end ribs 23 and the intermediate ribs 24 are respectively interposed so as to be in parallel to one side of the plate-like body 21. The upper edges of the ribs 23, 24 are jointed to the upper plate like body 21 or 22, and the lower edges of the ribs 23, 24 are jointed to the lower plate-like body 22 or 21, whereby a plurality of particulate-containing gas passages 27 with both ends opened are formed. In one side of the plate-like body 22, the ribs 23, 24 are formed, and in the other side of the same plate-like body 22, the ribs 25, 26 are formed extending in the direction perpendicular to the ribs 23 24. The ribs 25, 26 are formed substantially in the same manner as the ribs 23, 24 and differ in that they are formed in the direction perpendicular to that of the ribs 23, 24. Thus, a plurality of particulate-containing gas passages 27 and clean gas passages 28 with both ends opened are alternately formed, where the direction of passages of the former is perpendicular to the passages of the later.\nIn the filter unit, both ends of the particulate-containing gas passages 27 are opened, but the openings at their one side of the passages 27 are directly or indirectly closed, and an exhaust gas from a diesel engine is introduced from the openings of the other side of the filter unit 20. Or, alternatively, the exhaust gas from the diesel engine can be introduced from the openings opened at the both ends simultaneously. The plate-like bodies 22 function as a filter wall, and particulates are trapped on the surface of the plate-like bodies 22 at the side of the particulate-containing gas passages 27, whereby a clean gas without particulates passes through the plate like bodies 22 and then to the outside of the filter unit via the clean gas passages 28. The trapped carbonaceous particulates in the exhaust gas are removed by burning when the plate-like bodies 22 are heated at appropriate time intervals. Thus, the filter unit 20 is refreshed.\nIn the prior art technique as described above, the carbon particulates collected are burnt off by heating the plate-like bodies 22 having a filtering function to the ignition temperature of the particulates to thereby remove them. However, the plate-like bodies 22 are exposed repeatedly to a high temperature when the particles are periodically burnt off. Accordingly, sintering of the plate-like bodies is further accelerated, whereby a distribution of pore size originally formed changes to thereby increase a pressure loss during the particulates trapping operation. This causes difficulty in maintaining a stable performance of the filter unit. Further, a damage by fusion of the plate-like bodies 22 may take place when they are subjected to a high temperature at the time of burning off the particulates. When such damage occurs, it is no more possible to trap the particulates.\nThe particulates of exhaust gas from the diesel engine contains not only carbon but also an unnegligible amount of non-combustible components (for instance, 1%-5% by weight to the total amount of the particulates), and these non-combustible components are also trapped by the filter unit. Furthermore, the exhaust gas contains components such as SOx and NOx, which corrode materials constituting conduit elements to produce non-combustible solid components which deposit on the filter unit. These non-combustible solid components can not be removed by burning and accumulate in the filter unit to cause clogging, whereby properties of the filter unit are deteriorated.\nJapanese Unexamined Patent Publication No. 225721/1984 discloses a dust removing apparatus for a hot dust-containing gas. The dust removing apparatus comprises a plurality of vertical filter tubes having both ends opened. A dust-containing gas is introduced from the upper portion of the filter tubes. Dusts hindered by the filter tube walls from passing therethrough are moved downwardly in the filter tubes and are collected in a dust hopper provided below the filter tubes, while a clean gas obtained by passing through the filter tube walls is discharged from the side walls of the filter unit. It is understood that the dust removing apparatus disclosed in this document is suitable for treating a large flow rate of gas containing noncombustible combustible dusts discharged from an apparatus such as a converter in an iron plant. However, the document does not suggest application of the apparatus to the treatment of an exhaust gas from a diesel engine.\nSupposing this kind of system for a compact filter unit, it is easy to manufacture the simple shaped filter tubes, but there is such a problem that an assembling process requiring much labor and time is needed in order to fix a number of thin filter tubes to support plates with spaces among the tubes and it is not easy to manufacture a filter unit having a large capacity.\nThe applicant of this application proposed (in Japanese Patent Application No. 232360/1987) a filter device for an exhaust gas from the diesel engine as shown in FIG. 10 in order to solve the above-mentioned problem.\nA typical example of such device will be described.\nA filter unit 33 having a rectangular solid form which is substantially the same as the filter unit shown in FIG. 9 is used.\nThe filter unit is secured in a casing 31 having openings at upper and lower parts and a side port by interposing sealing members 32 as shown in FIG. 10. In passages and respectively have one end closed and the 34 which vertically pass through the filter unit 33 (as shwon by solid line arrow marks in FIG. 10) and clean gas passages 35 (as shown by broken line arrow mark in FIG. 10) which cross the particulate-containing gas passages and respectively have one end closed and the other end opened, and each passages are defined by plate-like bodies 22 made of a gas-permeable porous material.\nAn inlet duct 37 is formed at the upper part of the casing 31 to introduce an exhaust gas from the diesel engine. An outlet conduit 38 is connected to the casing 31 at a side where the clean gas passages 35 are open. The outlet conduit 38 is provided with a throat portion 39 having a reduced diameter part, and the portions contiguous to the upstream and the downstream of the throat portion are gently expanded. A nozzle 40 for ejecting a pressurized gas, which is used for back washing, is provided near the downstream side of the throat portion 39 so as to open toward the upstream side.\nA particulate receiving section 41 is provided at the lower part of the casing 31. The particulate receiving section 41 is provided with a tray 42, a filter plate 43 with an electric resistance type heater 46, an ash component removing port 44 with a removable cover plate 47 (it is usually closed), and a gas duct 45.\nA part of the bottom of the tray 42 is perforated and the filter plate 43 is fitted to that part so that the tray 42 and the filter plate 43 surround as a whole all the lower open ends of the particulate-containing gas passages 34 of the filter unit 33. The ash component removing port 44 is open at the bottom of the tray 42, and the gas duct 45 is placed below the filter plate 43. The filter plate 43 is made of a gas-permeable porous material. The ratio of a filtering area of the filter plate 43 to that of the filter unit 33 is such that about 20% or lower portion, especially about 0.5%-5% portion, therefore, small portion of the total exhaust gas introduced from the inlet duct 37 passes through the filter plate 43, and the remaining portion of the exhaust gas passes through the plate-like bodies 22 of the filter unit 33 to be flown to the outlet conduit 38.\nThe exhaust gas from the diesel engine is introduced from the inlet duct 37 via the upper open ends into the particulate-containing gas passages 34 of the filter unit 33. The most part of the exhaust gas passes through the plate-like bodies 22, the clean gas passage 35 and is flown to the outlet conduit 38. However, particulates in the exhaust gas are trapped by the plate-like bodies 22 and deposit on the inner surfaces of the particulate-containing gas passages 34. A part of the particulates flow to the particulate receiving section 41 together with a part of the exhaust gas through the lower open end of the particulate-containing gas passages 34. The part of the exhaust gas passes through the filter plate 43 to the gas duct 45. However, the particulates in the part of the exhaust gas can not pass through the filter plate 43 and deposit on the inner surface of the filter plate 43.\nAfter a continuation of the particulate trapping operation for an appropriate period, a short time back washing operation is carried out. In the back washing operation, a pressurized gas, especially pressurized air is ejected from the nozzle 40 for a short time such as about 0.1 sec-1 sec. The ejected gas induces the gas around the nozzle 40, and the gas of several times as much as the original ejected gas is flown in a pulse form to the clean gas passages 35. The pulse gas flow is flown to the particulate-containing gas passages 34 through the plate-like bodies 22. Then, the particulates accumulated on the inner surface of the particulate-containing gas passages 34 are peeled off. A part of the particulates drifts in the particulate-containing gas passages 34, however, the most part drops into the particulate receiving section 41. In the particulate receiving section 41, a stream of gas is produced to pass through the filter plate 43, and being carried with the gas flow, the most part of particulates deposit and accumulate on the inner surface of the filter plate 43.\nThus, the particulates trapped on the inner surfaces of the particulate-containing gas passages 34 during the particulate trapping operation are moved onto the inner surface of the filter plate 43 by the back washing operation to thereby refresh the filtering function of the filter unit 33. The particulates on the filter plate 43 are burnt and removed by heating the electric resistance type heater 46.\nAfter a relatively long term use of the apparatus, there occurs accumulation of non-combustible residue or ash. In this case, the cover plate 47 is opened to remove the residue or ash. Alternatively, they may be forcibly removed by a suitable scraping means.\nHowever, the above-mentioned filter device had a disadvantage as follows. When the back washing gas was introduced in a pulse form into the clean gas passages 35 repeatedly, a high gas pressure head such as about 0.3-0.4 kg/cm.sup.2 is instantaneously applied to the inside of the clean gas passages 35. Each application of the pressure head induces a stress to the plate-like bodies 22 which separate the clean gas passages 35 from the particulate-containing gas passages 34 in the filter unit 33 and finally damages the plate-like bodies 22.\nFurther, the conventional filter unit as shown in FIG. 9 had a complicated structure which requires much labor and long manufacturing time and the yield rate is low."}
{"text": "The present invention relates to mowers comprising bottom-driven discs, more particularly those where the discs are driven by means of bevel couplings. In this type of mower the bevel couplings are housed either in a single housing extending beneath all the discs and containing a bath of lubricant, or in housings where the lubricant is in danger of escaping when the transmission shaft extending through each disc housing is withdrawn therefrom.\nAll these machines possess one or more of the following drawbacks:\nLengthy and delicate assembly and dismantlement both for the manufacturer and especially for the user,\nNecessity of adjustment of the bevel coupling each time the transmission member passing through the disc housing must be re-fitted after repair,\nRisk of oil losses, especially during repairs,\nRisk of penetration of dirt into the bevel couplings and their bearings,\nRisk of incorrect oil level; it is known that even a slight upward or downward variation of oil level involves grave consequences in view of the high rotation speed of the bevel couplings driving the discs,\nRisk of poor lubrication of the end bevel couplings when the cutter bar is working in an inclined position, for example along a slope, and risk of deterioration of a large number of transmission members in the case of an accident to any one of them.\nThe purpose of the present invention is to remedy all these drawbacks."}
{"text": "The present invention relates to a 4-acyloxyquinoline derivative and an insecticide containing it as an effective ingredient.\nA variety of compounds having insecticidal or acaricidal activities have hitherto been developed. However, it is desired to develop a new insecticide having an excellent activity because of the problem such as the recent increase of insecticidal resistance."}
{"text": "The present invention relates to a group of new compounds of the general formula I:\nNu-O-Fa\nwherein O represents an oxygen, Nu is a nucleoside or nucleoside analogue, and Fa is an acyl group of a mono-unsaturated C18 or C20 fatty acid. The invention also concerns anti-viral pharmaceutical and veterinary compositions comprising a compound of formula I alone or in combination with a pharmaceutically acceptable carrier. A further part of this invention is a method for the treatment of a human or animal patient suffering from a viral infection and for reducing the infectious load by administering a compound of formula I.\nA large number of serious diseases, such as AIDS, hepatitis B, herpes and gynecological cancer, as a late result of papilloma warts, are caused by viral infections.\nViruses are small infectious agents which are incapable of independent replication and thus are dependent on a host cell to replicate. The genetic material of the virus is either RNA or DNA.\nWhen infecting an organism, the virus attaches to a specific host cell. The virus penetrates the cytoplasmic membrane after attachment and the viral genome is released from the virus particle. The viral genome is usually transported to the cell nucleus where new viral genomes are replicated. New viral protein is synthesized in the cytoplasm and new particles are formed either close to the cytoplasmic or nuclear membrane.\nSome viruses have genomic material which is directly (DNA virus) or indirectly (Reverse transcription of RNA, retrovirus) incorporated in the host cell genomes.\nExtracellular viruses are neutralized by circulating antibodies and the cellular immune apparatus may attack and remove infected cells. Viruses within the infected cells escape immune surveillance if viral antigens are not exposed on the surface of the cells.\nThe immune attack on infected organs contributes to disease by a mechanism commonly called virus induced immuno-pathology.\nThe mechanisms underlying some of the more important viral diseases differ.\nWhen suffering from a HIV infection, the patient\"\"s T helper cells are invaded and destroyed. This leads to a immunodeficiency condition, which makes the patient very succeptible even to infections which normally are conquered by the immune system without any harmful effects for the patient.\nThe Hepatitis B virus invades the liver cells, and the patient may become very ill when the immune system tries to rid the body of these infected cells. If the infection is not conquered by the immune system at an early stage, the result will be chronic hepatitis. The patient will thus be infectious throughout his life. For a group of patients the chronic hepatitis will develop into cirrhosis or cancer of the liver.\nIn herpes simplex infections, the virus enters the epidermal cells originally. The herpes simplex virus travels up to a nerve center where it lies latent to break out at intervals. Although not life threatening in most cases, a herpes infection is painful and the patient will be infectious every time an outbreak occurs.\nIn the papilloma virus, notably in the genital tract of women, the viral genome is located in the nucleus of epithelial cells, but not integrated in the cell chromosomes. This is a persistent condition and with some tumor promoting strains an integration finally occurs leading to a malignant development. The viral genome in this case has a decisive initiating effect in the process leading to cancer.\nIf the immune system manages to rid the body of the virus at an early stage, this leads to a life long immunity. On the other hand, if the virus is too aggressive and avoids the immune apparatus, no immunity is achieved and a continuous infectious state is the result.\nAs a result of the different mechanisms, the therapeutic strategy would be different for these conditions.\nThe ultimate goal in the treatment of HIV/AIDS would be to free the patient from the infectious virus. This seems to be remote at the present stage. However, much can be obtained by improving the general condition of the patient. A reduction of the virus load would increase the length of the symptom free period and reduce the infectiousness, which is of utmost importance in regard to the epidemiological situation. All currently used anti-viral agents have toxic side effects, which presently makes a sufficiently aggressive treatment impossible.\nIt is assumed that there are between 250 and 300 million carriers of hepatitis B world wide. It is known that a great number of these are going to develop hepatomas or liver insufficiencies due to the infections. Promising results in the treatment of the carrier state have been obtained in recent years by induction of an immune response with interferon. Therapies reducing the virus load are important in this regimen as efficient treatment of acute hepatitis B would reduce the number developing into a carrier state. The recently identified hepatitis C virus causes a very great number of cases with hepatitis whereof a large number develop into carriers. Preliminary studies seem to indicate that the carrier state may be broken by similar therapeutic regimens as for hepatitis B. Herpes simplex 1 and 2 frequently infect humans causing a carrier state with recurrencies of local infections. Generalized infections including encephalitis are rare but a catastrophy for the patient. There is a great individual variation in the frequency of local infections. For those patients who are affected either genitally or facially this constitutes a serious health problem physically, mentally and socially. None of the therapeutic regimens developed so far cures the latent infections of cells in the central nervous system. The therapeutic goal is thus to minimize the clinical manifestations of recurrencies both as to symptoms and duration.\nThe prevalence of genital papilloma virus infections has increased dramatically during the 1980s. It is now established that some genotypes are oncogenic, that is they initiate changes in the cell which after a latency period develop into cancer. Papilloma virus of the genital tract give long standing infections. The factors causing malignant transformation of the lesions are not well understood, but the immune system is assumed to be of importance. It is thought that the lesions showing progression during the months and years are those giving rise to cancer. The genital papillomas called condylomas are at present treated by physical means such as surgical removal, necrotizing means, fluid nitrogen or the like. Genital warts are at the onset benign tumors with altered enzyme patterns affecting among other things the metabolism of nucleoside analogues. Nucleoside prodrugs affect the episomal proliferation of papilloma virus thereby inducing regression of the warts.\nProphylactic vaccination has been very successful in acute infections such as polio, measles, mumps etc., but no effective vaccination has been developed for many of the other serious viral infections.\nEven though there have been intensive efforts to produce effective anti-viral chemotherapeutica during the last decades, no satisfactory medical treatment can be offered for most viral diseases today. The efforts have been especially great since the appearance of the HIV and related viral infections, which are spreading throughout the world at an alarming rate, yet the effects obtained with agents such as azidothymidine (AZT) and acyclovir (ACV) in AIDS and herpes can only be characterized as partially successful. These most promising anti-viral agents are thus derivatives of naturally occurring nucleosides, which have been modified either in the base or sugar moiety. They have, however, not had the therapeutic potential hoped for, as they bring forward serious side-effects in some patients or show little or no effect in others. Further, treatment with these agents is extremely expensive. For these reasons only patients suffering from the very serious viral infections such as AIDS receive such treatment. Patients suffering from the less serious, but also very painful viral infections are often left without any treatment to let the infection take its own course.\nThe untreated patient carries a great infectious load and constitutes a risk to his fellow human beings. If he is treated with an anti-viral agent, the aim is to reduce the infectious load so as to enable the body\"\"s immune system to conquer the infection. A further aim is to reduce the contagiousness and thus the number of new patients and carriers.\nThus the need for compounds having a better therapeutic index is obvious.\nThe need is especially great in chronic or recurrent viral infections with a dangerous acute phase or long term ill effects on health or well being, such as AIDS, hepatitis B and C, infections of the herpes group and papilloma viral infections. Similarily, there is also a need for anti-viral agents usable in the treatment of animals suffering from viral diseases.\nIn order to improve the effect, there have been developed derivatives of the nuclosides which are either modified in the base or sugar moiety. Especially fatty acid esters of the nucleosides or of analogues of nucleosides have been developed, in order to improve the lipophility and achieve a better membrane passage.\nThus there are known from EP 56265 (Lxc3x6bering et al.) esters of arabino-furanosyl-thymine (Ara T) with saturated acids having 1-17 C-atoms.\nFrom PCT/WO90/00555 (Hostetler et al.) there are known lipid derivatives linked, epecially through a phosphate group, to the 5xe2x80x2-position of the pentose group of a nucleoside. The purpose of this derivatization was to make the nucleosides more lipophilic so that they could be included into liposomes, which are preferentially taken up by macrophages and monocytes, cells which are found to harbor the HIV virus. It is stated that a targetting effect is thereby achieved.\nFrom EP 393920 there are known unsaturated, preferably polyunsaturated, C16 or higher fatty acid esters or amides of nucleoside or nucleoside analogues. It is stated that the fatty acid portion of these molecules preferably are made up of poly-unsaturated fatty acids, such as xcex3-linolenic or linoleic acid.\nIt has now surprisingly been found that a selected group of fatty acid, esters of anti-viral nucleosides or nucleoside analogues, wherein the fatty acid is a mono-unsaturated C18 or C20 acid gives a much improved effect. Also the compounds according to the present invention seem to be relatively nontoxic as judged by growth experiments with young mice and also have a low cytotoxicity according to observations made in tissue culture experiments.\nAlthough it is known that both nucleosides and nucleoside analogues, by themselves, and also some unsaturated fatty acids, by themselves, exhibit anti-viral effects, the magnitude of the effects achieved with the compounds according to this invention indicates that this is not an additive, but rather a synergistic activity which is special for the compounds of formula I.\nThe mechanism behind these effects is at present not known, but they are of an order of magnitude better than the closest related compounds of the prior art. Thus it is not considered likely that they arise only due to a membrane effect or targetting effect.\nFurther it is also clear, as will appear from the biological examples included herein, that effects are achieved with these compounds in systems where no effects may be achieved with the mother nucleoside compound.\nThe compounds of this invention can be characterized by the general formula I:\nNu-O-Fa\nwherein O is an oxygen, Nu is a nucleoside or nucleoside analogue and Fa is an acyl group of a mono-unsaturated C18 or C20 fatty acid.\nNucleosides are molecules comprising a heterocyclic base, such as cytosine, uracil, adenine or guanine, linked to a ribose unit. In nucleoside analogues, either the base or the ribose unit has been modified. For example, the ribose unit may be replaced by another sugar unit or by a noncyclic chain.\nThe fatty acid is esterified with a hydroxyl group in the 5-position of the sugar moiety of the nucleoside or with a hydroxyl group on the non-cyclic group of the nucleoside analogue.\nThe nucleoside or nucleoside analogues which may be chosen as Nu in the compounds of formula I may preferably be represented by the formula II:\nSxe2x80x94Bxe2x80x83xe2x80x83(II)\nwherein S is either a mono-saccharide-derivative selected from: 1-xcex2-D-arabinofuranose or 2,3-di-deoxy-3-azido-1-xcex2-D-ribofuranose, or selected from the group of 2-hydroxy-ethoxy-methyl, 4-hydroxy-3-(hydroxymethyl)-butyl, 2-hydroxy-1-(hydroxy-methyl)-ethoxy-methyl or 2,3-di-hydroxy-propoxy; and B is a nitrogen base selected from adenine, guanine, cytosine, uracil, thymine or a thymine derivative of the following formula: \nWherein X is deuterium or fluorine.\nExamples of these nucleosides or nucleoside analogues are: \nExamples of the group R in different nucleoside analogues appear below: \nThere exist several systems for denomination of the position of double bonds in fatty acids. In the present application the xcfx89-system is used, wherein the position of the double bond in the unsaturated fatty acids is counted from the terminal methyl group. Eicosenic acid (C20:1 xcfx89-9), for example, has 20 carbon atoms in the chain, and the double bond is found between carbon atom 9 and 10 counting from the end of the chain.\nThe fatty acids which are reacted with the nucleosides or nucleoside analogues to form the esters according to this invention with the pronounced activity are the xcfx89-9 C18 or C20 monounsaturated fatty acids. Further, even though the effect observed differs somewhat between acids of the same chain length when the double bond is in cis or trans configuration, both show a strong activity. Also the position of the double bond is of importance, as it appears that especially the esters with a C18 or C20 fatty acid having the unsaturation in xcfx89-9 position will have the surprisingly elevated activity.\nThe C18 or C20 xcfx89-9 fatty acids, which when bound to the nucleosides give the surprisingly elevated effect, are the following: oleic acid (C18:1, xcfx89-9, cis), elaidic acid (C18:1, xcfx89-9, trans) eicosenic acid, (c20:1, xcfx89-9, cis) and (C20:1, xcfx89-9, trans)\nPreferred representatives of the compounds according to this invention are listed below.: Ara A-oleic acid, Ara A-elaidic acid, Ara A-eicosenic acid, cis, Ara A-eicosenic acid, trans, Ara T-oleic acid, Ara T-elaidic acid, Ara T-eicosenic acid, cis, Ara T-eicosenic acid, trans, ACV-oleic acid, ACV-elaidic acid, ACV-eicosenic acid, cis, ACV-eicosenic acid, trans, AZT-oleic acid, AZT-elaidic acid, AZT-eicosenic acid, cis, AZT-eicosenic acid, trans. Their formulas will appear from FIG. 4.\nThe compounds according to this invention exhibit anti-viral effects, and the present invention thus includes pharmaceutical or veterinary compositions comprising at least one compound of formula I alone or in combination with a pharmaceutically acceptable carrier or excipient. In the remainder of the text and in the claims, a pharmaceutical composition will be used for compositions usable in the treatment of both human and animal patients.\nFurther, it appears that certain of the monounsaturated fatty acid nucleosides or nucleoside analogues will be especially suitable for the treatment of certain viral infections. Thus it appears that the fatty acid derivatives of AZT are especially useful for the treatment of AIDS.\nSimilarly, it appears that the fatty acid derivatives of Ara T and Ara A are especially suitable for the treatment of hepatitis B. It appears also that the Ara T esters will be efficient in agents suitable for treatment of papilloma viral infections.\nFurther, it appears that the fatty acid esters according to this invention of ACV are especially suitable for the treatment of herpes infections.\nAs mentioned, the production of the necessary immune response in order to conquer a viral infection, such as hepatitis, can be induced in some cases by the co-administration of interferon.\nDepending on which viral infection is to be treated and at what stage the infection is, or if the patient is a human being or an animal, both a systemic and a local administration of the compounds can take place.\nFor local administration, the compounds can be formulated as known in the art for administration to the skin or mucosa in any suitable form.\nWhen administered topically the compounds of formula I may be formulated as an ointment, cream, gel, tincture, spray, lotion or the like containing the compounds of formula I in admixture with inert, solid or liquid carriers which are usual in topical preparations. It is especially suitable to use a formulation which protects the active ingredient against oxidation or degradation.\nThe pharmaceutical preparations comprising the compounds of formula I may also be administered systemically, either enterally or parenterally.\nWhen administered enterally, the compounds of formula I may be formulated e.g. as soft or hard gelatine capsules, tablets, granules, grains or powders, dragees, syrups, suspensions or solutions.\nWhen administered parenterally, preparations of the compounds of formula I as injection or infusion solutions, suspensions or emulsions are suitable.\nThe preparations can contain inert or pharmacodynamically active additives. Tablets or granulates e.g. can contain a series of binding agents, filler materials, carrier substances or diluents. Liquid preparations may be present, for example, in the form of a sterile solution. Capsules can contain a filler material or thickening agent in addition to the active ingredient. Furthermore, flavor-improving additives as well as the substances usually used as preserving, stabilizing, moisture-retaining and emulsifying agents, salts for varying the osmotic pressure, buffers and other additives may also be present.\nThe dosages in which the preparations according to this invention are administered will vary according to the mode of use and the route of use, as well as to the requirements of the patient. In general a daily dosage for a systemic therapy for an adult average patient will be about 0.1-100 mg/kg body weight/day, preferably 1-20 mg/kg/day. For topical administration, the suitable ointment can contain from 0.1-10% by weight of the pharmaceutical formulation, especially 0.5-5% by weight.\nIf desired, the pharmaceutical preparation of the compound of formula I can contain an antioxidant, e.g. tocopherol, N-methyl-tocopheramine, butylated hydroxyanisole, ascorbic acid or butylated hydroxytoluene.\nThe invention further discloses a method for the treatment of viral infections, which comprises administering at least one compound of formula I to a human or animal patient in need of such treatment.\nFurther, the invention also comprises a method for the treatment of a patient in need of such treatment with a combination of a compound of formula I and an interferon."}
{"text": "This invention relates to curing of phenolic novolac resins at high temperatures. More particularly, this invention relates to the use of (lower) alkoxylated melamine-formaldehyde resins and (lower) alkoxylated benzoguanamine-formaldehyde resins as curing or hardening agents for novolacs at high temperature. The curing agents of this invention are also simply referred to herein as triazine hardeners.\nIn, refractory applications, there is a need for a binder comprising a phenolic novolac and a curing agent wherein the binder remains uncured until high temperatures are attained. In the binders of this invention at least 50% of the cure and preferably at least 80% of the cure is accomplished above 392.degree. F. (200.degree. C.). This minimizes hardening and solidification of the binder in heated equipment and provides improved flow properties and prolonged workability for the binder.\nThere is also a need for such binder compositions which can be delivered to hot surfaces under pressure wherein the compositions further contain a particulate refractory and water wherein the hardener is not leached out from intimate contact with the other ingredients of the composition. This prevents or minimizes stratification on drying and curing of the composition.\nPitch as well as phenolic resins are used as binders for aggregates to prepare compositions for refractory tap hole application for blast furnaces. Tap hole mix users need a combination of early green strength and stability at 150.degree. F. (65.degree. C.) to 250.degree. F. (121.degree. C.) for up to several hours. Novolac with hexamethylenetetramine (hexa) has been used in such compositions. Also, hexa has been used as a hardener for novolac in the manufacture of electrodes. However, novolac/hexa binders lead to early cure whereas novolac without hexa or other hardener do not cure. Resoles are also unsatisfactory in such compositions because they cure at too low a temperature.\nMelamine resins have also been mentioned as curing agents for novolacs or as thermosetting resins along with resoles, etc. in refractory compositions. However, the melamine resins used in this invention have advantageous properties in their use.\nIt is known that novolacs containing acids or thermally produced latent acids strongly catalyze the cure of melamine resins. However, the acid lowers the normal curing temperature of the novolac.\nApparatus and procedures exist for the robotic spraying of refractory material in a slurry of about 20% water for maintenance and repair of heated refractory surfaces. Illustrative of such a system is the TUNCAST Tundish Spray System of Foseco Inc."}
{"text": "In gas turbine engines, the operational clearances between static engine structures and the tips of rotating blades impact performance aspects such as the thermodynamic efficiency and fuel burn of the engine. It is desirable to reduce the operational clearances, while at the same time avoiding contact between the rotating blade tips and the static structure. The performance improvement value of as little as a tenth of a millimeter in clearance reduction is significant. In response to thermal excursions, the lengths of the blade tips typically vary at a different rate than that at which the static structures expand or contract. Variation in position or concentricity leads to a need for greater clearances and suboptimal performance. For example, without excess clearance, the blade tips may contact the static structure negatively impacting engine performance. Excess clearance between the blade tips and the static structure may also reduce engine performance. When engine designs result in higher temperatures such as to achieve greater efficiency, maintaining optimal clearance without excess cooling air becomes more challenging.\nAccordingly, it is desirable to provide systems that effectively retain and control the position of the static structure relative to the blade tips. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description of the invention and the appended claims, taken in conjunction with the accompanying drawings and this background of the invention."}
{"text": "The invention relates to the field of cutlery and it has as an object a relief of a blade of a cutting tool, in particular of a knife.\nThere are known profiles which are arranged on the blades in order to level their cutting loss due to use and wear, which profiles are constituted by microtoothing which is composed of a plurality of incurved surfaces having parallel generatrixes, called microbevels, defining with one of the side faces of the blade a plurality of respective microcutters.\nThe inconvenience of the conventional microtoothing, formed by a plurality of short microbevels arranged along the blade, is based on the fact that the resulting short microcutters form a toothing which induces a shredding of the material to be cut up, where the tool therefore practically no longer cuts and, finally, being of a delicate use since the material, in particular meat, is soft.\nIn order to remedy this inconvenience, the efforts of the innovators has been directed to the refining of a microtoothing composed of short microbevels and of long microbevels, succeeding each other alternatingly, wherein the first shred the matter while the second cut the matter. One could in particular refer to the document EP-A-0220362 which describes a tool blade comprising such a microtoothing. In addition, the blade carries along its other face a bevel for reinforcing its cutting characteristic.\nThe result obtained is satisfying with regard to the improvement of the performance of the tool for cutting the matter but retains nevertheless another general drawback of microtoothed tools, namely to make their utilization delicate based on the presence of the microbevels which induce a tendency of the blade to incline itself by turning its blade edge toward the microbevelled face.\nIn addition, tools are known which carry a microtoothing composed of a plurality of short microcutters, successively formed by the short microbevels arranged respectively oppositely on each one of the side faces of the blade, in order to confer to the tool a natural stability during the time of its utilization.\nThe use of this type of tool, which is of relatively ancient origin, has been rapidly abandoned because of its high manufacturing cost, brought about by the difficulty of placing oppositely the microbevels arranged on the one and the other face of the blade, resulting in an unsatisfactory cut based on the cutter absence on the blade."}
{"text": "Surface measurement systems are used in a variety of applications to generate three dimensional surface data of objects. Such systems are employed at various stages in the fabrication and assembly of complex objects across a variety of industries to ensure that the shape and size of the objects meet strict manufacturing tolerances.\nInterferometric surface measurement systems have been developed which permit measurements of the surface of an object without physical contact. Coherent optical sources are used to generate a fringe pattern on the surface of the object and a camera acquires images of the fringes on the surface for analysis. In some systems, a diffraction grating is positioned in the path of a laser beam to generate multiple coherent laser beams at various angles to the original beam path. A focusing objective and spatial filter are used to isolate the desired diffracted beam pair. One or more additional diffraction gratings are utilized to project at least one additional set of fringes onto the object surface. This multiplexing of different gratings into the beam path poses many challenges. Moving different gratings into the beam path and shifting each grating to implement phase shifts generally requires multiple mechanical components that add weight, size, complexity and cost to the system. The frequent movement of components affects the stability and therefore the accuracy of the measurement data. Moreover, measuring the displacement of a diffraction grating during the phase shift process with sufficient precision and accuracy can require expensive measurement components such as capacitance gauges.\nOther system components can limit the applications for the system. For example, the focusing objective and spatial filter are used for multiple gratings and, therefore, their optical parameters are not optimal for the individual gratings. Moreover, the depth of field of the camera can limit the maximum spatial frequency of the projected fringe pattern, thereby limiting the measurement resolution.\nNoise sources also typically limit the measurement data. When laser light is scattered from a surface, a high-contrast, granular speckle pattern is typically observed. Speckle results in part from the roughness of the object surface. In particular, the microscopic roughness of the surface contributes randomly phased contributions of the scattered laser light. These contributions interfere with one another to produce complex intensity variations across the surface of the object as viewed from a distance. Speckle introduces fine-scale intensity fluctuations (i.e., intensity noise) in the observed fringe pattern on the object surface. Shot noise contributions from the individual detectors in the camera can further limit the accuracy of measurement data."}
{"text": "The present invention relates to a semiconductor device, particularly to a multi-chips module (MCM) type semiconductor device adapted to be mounted on a board and to be electrically connected to the board.\nJP-A-11-220077 discloses that a coefficient of thermal expansion and so forth is adjusted to restrain a crack of a semiconductor element and/or an under-fill in a flip-chip type semiconductor device. JP-A-2000-40775 discloses a shape of an oblique surface of the under-fill is adjusted to restrain the crack of the semiconductor element."}
{"text": "1. Technical Field\nThe present disclosure relates to data networks providing access to data storage in general and, in particular, to optimization of power consumption by network nodes operating as data storage.\n2. Description of the Related Art\nMedia, such as audio recordings and films, can nowadays be effectively archived in digital format, and it is used by private persons and by large organizations. In the case of large organizations the number of records (i.e., individual audio or video recordings) can be very large. Additionally, the size of digitized media files, especially video files converted to high quality digital format, is large, and in consequence archives of media files require large memory space. The digitized media can be stored on various media, e.g., CD-ROM, DVD, Blue Ray, HD DVD, hard disks, magnetic tapes, and others.\nFor large organizations using hard disks for storing media data is very useful as it does not require any manual or mechanical operations to be performed in order to access the files stored on that drive, and as a result it can be accessed from a remote location (e.g., a branch of this organization) via a network (e.g., intranet or Internet).\nThe problem faced by some organizations is that the number of media files is growing very fast, and the old files cannot be simply deleted because, for example, they are classified as part of national heritage and as such must be archived indefinitely. This results in a constantly expanding archive. The problem of available memory space can be solved by connecting additional hard disk drives to the existing ones.\nWhile the solution with a plurality of hard disk drives may provide adequate space for storing the media files, it also requires a substantial amount of electricity to run. Power consumption of a storage system comprised of 10,000 hard disks, each consuming 25 W is 250 kW. By reducing power consumption of such a system, significant financial benefits can be gained, and additionally its impact on the environment will be reduced."}
{"text": "1. Field of the Invention\nThe present invention generally relates to gift cards, and, more particularly, to systems and methods for gift card linking.\n2. Description of the Related Art\nGift cards are very popular, and are often among the most requested type of gift. Gift cards provide advantages for everyone involved in the gift-giving process. For example, the gift card giver does not have to spend time determining the gift that the gift receiver wants, needs, or has, but can instead give a gift card to one of the gift receiver's preferred merchants. The gift card receiver can select what he or she wants, and does not have to worry about returns or exchanges. The merchants often see an increase in sales, as the gift card receivers often spend more than the amount of the gift card."}
{"text": "An automobile has an openable/closable mechanism including a box and a lid such as a glove box. A damper is provided to the edge of a box the glove box. When a lid of the glove box is closed, the damper functions to suppress the collision shook between the lid and the box with the elastic force by abutting on the lid.\nThe damping part described in Patent Document 1 includes a casing having a bottomed cylindrical shape, a retainer having a bottomed cylindrical shape and housed in the casing, and a spring abutting on the bottom of the casing and the bottom of the retainer to bias the retainer in a direction that the retainer protrudes out from an opening portion of the casing. A retaining protrusion provided on a lateral face of the retainer is engaged into an engagement window provided on a lateral face of the casing, and thereby retaining the retainer the biasing force of the spring."}
{"text": "The present invention relates generally to the field of tactile feedback technology, and more particularly to creating a visual representation of tactile feedback.\nHaptic technology, or haptics, is a tactile feedback technology that recreates the sense of touch by applying forces, vibrations, or motions to a user through actuators (e.g., motors responsible for movement) incorporated in a computing device. Haptics manipulate the actuator using defined waveforms to produce a range of effects which may be perceived uniquely by an individual touching the computing device. Initially, haptics were only capable of providing strong feedback with a limited range of sensations to an entire device or providing limited localization of haptic feedback to a specific position (e.g., display instead of the body of the device) with an expanded effects, such as frequency range, response time, and intensity. As the technology continues to evolve, haptics currently incorporated in computing devices now provide the means to deliver both touch-coordinate specific responses and customizable haptic effects (e.g., virtual keyboard scrolling list) to the user.\nWhen sensing physical interaction between a user and a computing device is necessary, haptic computing devices incorporate tactile sensors that measure information arising from a physical interaction with the environment. One of the most common implementations of tactile sensors are touchscreen devices, such as those utilized in mobile devices and computing. The touchscreen device allows a user to interact directly with what is displayed (e.g., images, icons, virtual keyboards, applications, etc.) by touching the screen with a specialized stylus and/or one or more fingers rather than utilizing an intermediate device, such as a computer mouse or a touchpad. The mechanical stimulation provided by the haptic feedback may then be utilized to assist in the creation of virtual objects in a computer simulation (e.g., developed model representing characteristics, behaviors, and functions of a physical or abstract system over time). The continued exchanges between the user and computer simulation via the touchscreen may then provide additional feedback (e.g., changes and updates to images, text responses, change in displayed information, etc.), which allows for confirmation of actions taken by a user and for additional interactions to occur based on received responses.\nDigital images created and displayed by a computing device include metadata. Metadata provides information regarding the data content of the actual digital image. For example, digital images may include metadata that describes how large the picture is, the color depth, the image resolution, when the image was created, and other relevant data. In addition to the aforementioned metadata, additional metadata (e.g., GPS coordinates, compression ratios, software, etc.) and object tags (e.g., non-hierarchical keyword or term assigned to a digital image describing an item) may be changed and assigned to the digital image by a user and through image processing software which includes incorporating tactile effects. For example, a user or image processing software may identify objects and the boundaries associated with the object within a digital image. Once the boundaries of an object are defined, additional tactile effects may then be added to the metadata (e.g., tactile metadata), enhancing the viewing experience of future users."}
{"text": "Modern communications systems permit callers and callees, who are often separated by great distances, to freely dialogue with one another. However, when a caller initiates a call (or session) to a callee, very little information about the callee may be available to the caller prior to establishing a session. Similarly, when a callee receives a call from a caller, very little information about the caller may be available to the callee until a session has been established. In some instances, additional information about a caller or callee may facilitate dialogue between the two parties. To date, very little is available to provide such information in a communication system. A need, therefore, exists in the industry to address this deficiency."}
{"text": "Mobile communication devices such as smartphones, tablet computers, wearable devices, etc. are widely used by consumers for many things other than voice communications. Specifically, these devices are designed so that consumers can install various applications, which may perform many different tasks. For example, various applications may be used for things such as shopping lists, text messaging, social network communications and sharing, research, weather reporting, news, and so forth.\nApplications such as this often rely on interactions with remote servers or data stores, and may further rely on wireless data connectivity to communicate with such remote servers or data stores. For example, a weather application may communicate with a remote server in order to obtain weather information. Similarly, a list application may store lists on a remote server or database so that the lists can also be accessed on multiple devices.\nIn many situations, data connectivity is available to mobile devices through cellular communications networks, which are provided by cellular communication providers. Cellular communication providers maintain large infrastructures for providing wireless data communications for mobile devices. A cellular services infrastructure includes geographically distributed base stations as well as centralized communication services. In addition, many cellular communication providers also provide various consumer applications for installation and use on the mobile devices of their customers. For example, a cellular communication provider may provide an application that allows customers to manage their cellular service accounts. Customers may also obtain, install, and use many other applications, many of which may rely on the availability of cellular data communications."}
{"text": "Polysaccharides constitute the bulk of the plant cell walls and have been traditionally classified into three categories: cellulose, hemicellulose, and pectin. Fry, (1988) The growing plant cell wall: Chemical and metabolic analysis, New York: Longman Scientific & Technical. Whereas cellulose is made at the plasma membrane and directly laid down into the cell wall, hemicellulosic and pectic polymers are first made in the Golgi apparatus and then exported to the cell wall by exocytosis. Ray, et al., (1976), Ber. Deutsch. Bot. Ges. Bd. 89:121-146. The variety of chemical linkages in the pectic and hemicellulosic polysaccharides indicates that there must be tens of polysaccharide synthases in the Golgi apparatus. Darvill et al., (1980) The primary cell walls of flowering plants. In The Plant Cell (N. E. Tolbert, ed.), Vol. 1 in Series: The biochemistry of plants: A comprehensive treatise, eds. P. K. Stumpf and E. E. Conn (New York: Academic Press), pp. 91-162.\nEven though sugar and polysaccharide compositions of the plant cell walls have been well characterized, very limited progress has been made toward identification of the enzymes involved in polysaccharides formation, the reason being their labile nature and recalcitrance to solubilization by available detergents. Sporadic claims for the identification of cellulose synthase from plant sources were made over the years. Callaghan, and Benziman, (1984) Nature 311:165-167; Okuda, et al., (1993) Plant Physiol. 101:1131-1142. However, these claims were met with skepticism. Callaghan and Benziman, (1985), Nature 314:383-384; Delmer, et al., (1993), Plant Physiol. 103:307-308. It was only relatively recently that a putative gene for plant cellulose synthase (CesA) was cloned from the developing cotton fibers based on homology to the bacterial gene. Pear, et al., Proc. Natl. Acad. Sci. USA 93:12637-12642; Saxena, et al., (1990), Plant Molecular Biology 15:673-684; see also, WO 9818949; see also, Arioli, et al., (1998). Molecular analysis of cellulose biosynthesis in Arabidopsis. Science Washington D.C. Jan. 279, 717-720. A number of genes for cellulose synthase family were later isolated from other plant species based on sequence homology to the cotton gene (Richmond and Somerville (2000), The cellulose synthase superfamily, Plant Physiology 124:495-498.)\nCellulose, by virtue of its ability to form semicrystalline microfibrils, has a very high tensile strength which approaches that of some metals. Niklas, (1992), Plant Biomechanics: An engineering approach to plant form and function, The University of Chicago Press, p. 607. Bending strength of the culm of normal and brittle-culm mutants of barley has been found to be directly correlated with the concentration of cellulose in the cell wall. Kokubo, et al., (1989), Plant Physiology 91:876-882; Kokubo, et al., (1991) Plant Physiology 97:509-514.\nAlthough stalk composition contributes to numerous quality factors important in maize breeding, little is known in the art about the impact of cellulose levels on such agronomically important traits as stalk lodging, silage digestibility, or downstream processing. The present invention provides these and other advantages."}
{"text": "Projectors are used for presentations in school classrooms, meeting rooms and the like, as well as in private home theatres. As a light source, such projectors use a high-pressure discharge lamp that is nearly a point light source. In particular, high-pressure mercury lamps are used.\nA high-pressure mercury lamp encloses a halogen substance along with mercury, and includes an arc tube in which a pair of tungsten electrodes are disposed in mutual opposition. Light is radiated by the emission of an arc discharge between the electrodes. In a high-pressure mercury lamp, the halogen cycle is used to constrain the blackening effect produced during lighting. The blackening effect refers to the tungsten making up the electrodes evaporating and coming to adhere to the inside of the arc tube. The tungsten so evaporated is made to return to the electrodes through the halogen cycle, and comes to accumulate. As such, given a protrusion formed by such accumulation on the tip portions of both electrodes, the luminescent spots of the arc discharge spanning the inter-electrode distance are the two protrusions.\nIn recent years, demand for ever-increasing longevity in the high-pressure mercury lamp serving as the light source has grown alongside projector use. As a response to such demand, an effective method has been discovered in preserving the stability of the arc discharge by maintaining an appropriate shape in the protrusions formed on the tip portions of the pair of electrodes.\nOne proposed method for forming and maintaining an appropriate shape for said protrusions involves a high-pressure discharge lamp lighting device configured to switch between two or more frequencies for the AC current supplied to the high-pressure mercury lamp (see, for example, Patent Literature 1). Patent Literature 1 describes causing growth in the protrusions by applying suitable AC current in at least one frequency, and causing reduction in the protrusions by applying suitable AC current in other frequencies. The shape of the protrusions is thus preserved by repeatedly alternating between the frequencies to cause growth and reduction."}
{"text": "It is known that organic and polymeric materials with large delocalized .pi.-electron systems can exhibit nonlinear optical response, which in many cases is a much larger response than by inorganic substrates.\nIn addition, the properties of organic and polymeric materials can be varied to optimize other desirable properties, such as mechanical and thermoxidative stability and high laser damage threshold, with preservation of the electronic interactions responsible for nonlinear optical effects.\nThin films of organic or polymeric materials with large second order nonlinearities in combination with silicon-based electronic circuitry have potential as systems for laser modulation and deflection, information control in optical circuitry, and the like.\nOther novel processes occurring through third order nonlinearity such as degenerate four-wave mixing, whereby real-time processing of optical fields occurs, have potential utility in such diverse fields as optical communications and integrated circuit fabrication.\nOf particular importance for conjugated organic systems is the fact that the origin of the nonlinear effects is the polarization of the .pi.-electron cloud as opposed to displacement or rearrangement of nuclear coordinates found in inorganic materials.\nNonlinear optical properties of organic and polymeric materials was the subject of a symposium sponsored by the ACS division of Polymer Chemistry at the 18th meeting of the American Chemical Society, September 1982. Papers presented at the meeting are published in ACS Symposium Series 233, American Chemical Society, Washington, D.C. 1983.\nThe above recited publications are incorporated herein by reference.\nOf more specific interest with respect to the present invention embodiments is prior art relating to polymers with comb-like side chains. Eur. Polym. J., 18, 651(1982) describes liquid crystalline polymers of the smectic and nematic types with cyanobiphenyl groups in the side chain: ##STR2## where R is hydrogen or methyl, n is an integer of 2-11, and X is an oxy, alkylene or carbonyloxy divalent radical.\nSPIE Vol. 682, pages 56-64, Molecular and Polymeric Optoelectronic Materials: Fundamentals and Applications (presented at Aug. 21-22, 1986 meeting) describes liquid crystalline polymeric systems which are copolymers of a mesogenic monomer and a nonlinear optically responsive monomer.\nA disadvantage of liquid crystalline polymers which exhibit mesogenic side chain nonlinear optical response is an observed light scattering effect when the polymer is in the form of a solid phase optical medium, e.g., the polymer medium exhibits more than about 20 percent scattering of transmitted incident light. The light scattering is due to deviations from ideal molecular order which accommodate defects that are not optically clear.\nThere is continuing interest in the theory and practice of polymers which are characterized by comb-like side chain structures which can be oriented in an applied external field.\nThere is also an increasing research effort to develop new nonlinear optical organic systems for prospective novel phenomena and devices adapted for laser frequency conversion, information control in optical circuitry, light valves and optical switches. The potential utility of organic materials with large second order and third order nonlinearities for very high frequency application contrasts with bandwidth limitations of conventional inorganic electrooptic materials.\nAccordingly, it is an object of this invention to provide novel polymers with comb-like side chains.\nIt is another object of this invention to provide acrylic copolymers having side chains which exhibit nonlinear optical response.\nIt is a further object of this invention to provide optical light switch and light modulator devices with a transparent polymeric nonlinear optical component comprising an isotropic acrylic copolymer.\nOther objects and advantages of the present invention shall become apparent from the accompanying description and examples."}
{"text": "1. Field of the Invention\nThis invention relates to an image forming apparatus which has the function of detecting a jam of a recording material.\n2. Description of the Related Art\nFIG. 7 illustrates a sheet conveying path from a sheet-feeding unit to a sheet-discharging unit in a conventional image forming apparatus. In FIG. 7, there are shown sheet-feeding cassette 1 of a sheet-feeding unit, and transfer paper P, serving as paper accommodated within the sheet-feeding cassette 1. Sheet-feeding roller 2 is rotatably driven in response to a print start signal, whereby sheets of transfer paper P are individually taken out from within sheet-feeding cassette 1, and each sheet is fed to a transfer process portion of image forming unit 10 passing along a sheet path 6.\nIn this case, image forming unit 10 has an electrophotographic process mechanism. A drum-shaped electrophotographic photosensitive member (hereinafter termed a \"photosensitive drum\") 7 is rotatably driven in the clockwise direction indicated by the arrow at a predetermined circumferential speed (process speed).\nPrimary charging by charging unit (charging roller) 11, image exposure L by exposure means (not shown), such as a laser scanner or the like, and development by developing unit 12 are performed for rotating photosensitive drum 7, and a toner image corresponding to object image information is formed on the circumference of rotating photosensitive drum 7.\nThe toner image is sequentially transferred onto a sheet of transfer paper P fed from sheet-feeding cassette 1 at transfer nip portion 9 provided between photosensitive drum 7 and transfer unit (tranfer roller) 8.\nThe sheet of tranfer paper P passing through transfer nip portion 9 is then separated from the surface of rotating photosensitive drum 7, and is guided to fixing unit 17 passing along a sheet path 16. The surface of rotating photosensitive drum 7, after separation of the sheet, is cleaned by cleaning unit 13 to be repeatedly used for image formation.\nThe fixing unit 17 of the present case comprises a tensionless film-heating-type fixing unit disclosed, for example, in Japanese Patent Laid-open Application (Kokai) Nos. 4-44075-44083 (1992). FIG. 8 is an enlarged view of fixing unit 17.\nThe fixing unit 17 includes a heat-resistant film guide member 18 having a substantially semicircular cross section, a plate-like ceramic heater 19 having a low heat capacity disposed along the longitudinal direction at a central portion of the lower surface of film guide member 18, a cylindrical (endless) thin film (fixing film) 20 made of heat-resistant resin loosely fit around film guide member 18 having the ceramic heater, and a pressing roller 21 disposed below film guide member 18 so as to be in pressure contact with ceramic heater 19 across film 20 by always being pressed upwardly by upwardly-pressing spring 22. That is, ceramic heater 19 and pressing roller 21 are in pressure contact across film 20, to form a fixing nip portion N.\nCeramic heater 19 is heated by a current supplied from a current supply system (not shown), and is controlled at a predetermined fixing temperature by a temperature control system (not shown).\nPressing roller 21 is rotatably driven by a driving unit (not shown) in the counterclockwise direction indicated by the arrow at a predetermined circumferential speed. The cylindrical film 20 is rotatably driven around film guide member 18 by the rolling frictional force of pressing roller 21 in the clockwise direction indicated by the arrow at a predetermined circumferential speed while sliding along the surface of ceramic heater 19 in a state of firmly contacting the lower surface thereof at the fixing nip portion N.\nThe sheet of tranfer paper P conveyed to fixing unit 17 after image transfer is guided at entrance guide 23 to enter a space between the rotatably driven cylindrical film 20 and pressing roller 21 at the fixing nip portion N provided by the temperature-controlled ceramic heater 19 and the rotating pressing roller 21, and passes through the fixing nip portion N together with film 20 in a state of firmly contacting the lower surface of ceramic heater 19 across film 20.\nWhile the sheet of transfer paper P passes through the fixing nip portion N, the unfixed toner image on the sheet is heated across film 20 by the heat of ceramic heater 19. Thus, the image is fixed.\nThe sheet of transfer paper P passing through the fixing nip portion N is separated from the surface of film 20, and is output outside the apparatus after passing through exit guide 24, a sheet path 29 and a pair of sheet-discharging rollers 30.\nReferring again to FIG. 7, registration sensor S1 for detecting the leading edge of the sheet is disposed at a side upstream from transfer unit 8 in the conveying direction of the sheet. Discharged-sheet sensor $2 four detecting the discharge of the sheet is disposed at a side downstream from fixing unit 17 in the conveying direction of the sheet.\nAs shown in FIG. 8, the discharged-sheet sensor S2 of the present case includes a sensor lever 26 swingable around supporting shaft 25, and a photo-interrupter 27. Usually, sensor lever 26 has a substantially vertical posture by the function of gravity as indicated by the solid lines, with the upper-end portion thereof protruding in the sheet path and a lower-end shutter portion 26a thereof blocking the optical path of photo-interrupter 27.\nWhen the sheet is fed from the fixing nip portion N to exit guide 24, sensor lever 26 is pushed by the sheet to be rotated around supporting shaft 25 in the counterclockwise direction as indicated by the two-dot chain lines, whereby shutter portion 26a rotates to open the optical path of photo-interrupter 27.\nThus, the presence/absence of the sheet is detected by the opening/blocking of the optical path of photointerrupter 27. The registration sensor S1 has the same configuration as the discharged-sheet sensor S2.\nSignals from registration sensor S1 and discharged-sheet sensor S2 are input to a CPU (not shown) for controlling the main body of the image forming apparatus. Based on these input signals, control of the main body of the apparatus, such as timing control of exposure by a laser when exposure means for the rotating photosensitive drum 7 comprises, for example, a laser scanner, timing control of sheet feeding, control of the fixing sequence, jam detection and the like, is performed.\nA heat-roller-type fixing unit comprising a fixing roller, in which a halogen-lamp heater is provided within a roller having the shape of a hollow pipe, and a pressing roller in pressure contact with the fixing roller, a pressure fixing unit comprising a pair of pressing rollers, or the like may be used as the fixing unit 17. However, the film-heating-type fixing unit used in the present case is advantageous over other fixing units in that, for example, waiting time can be shortened (quick-start property) or electrical power consumption can be reduced, since a low-heat-capacity heating member, in which temperature rises rapidly, and a thin film can be used.\nIn the above-described image forming apparatus, if a paper jam occurs, the CPU for controlling the main body of the apparatus detects the jam from the relationship between the sheet-feeding timing and the signals from the sensors S1 and S2, stops the driving of the main body of the apparatus, and displays the occurrence of the jam.\nThe operator then opens an openable/closable cover of the main body of the apparatus, removes the jammed paper within the apparatus, and closes again the cover, whereby a door switch linked with the cover is closed. The CPU resets the jammed state by a signal from the door switch. By this reset, the main body of the apparatus executes pre-rotation driving as a mode of forcedly discharging paper, and returns to a state in which image formation can be executed.\nIn the above-described image forming apparatus, when the openable/closable cover of the main body of the apparatus is first opened after the driving of the apparatus has stopped based on the detection of a jam, and the cover is closed again carelessly or due to some other reason without removing the jammed paper, or when the jammed paper is insufficiently removed, and the cover is closed again without recognizing the presence of a remaining part of the jammed paper, the jammed state is reset based on a signal from the closed door switch, and the pre-rotation driving of the apparatus is executed, whereby the jammed state of the paper still remaining in the apparatus becomes worse, causing, for example, damage to the apparatus. Such problems also arise in the case of reset of the apparatus by the turning-on and turning-off of the power supply of the apparatus.\nThis is because when, for example, the paper produces an accordion jam, or a jam P' is produced by the paper becoming wound around the fixing member (comprising rotating film 20, pressing roller 21, the heat roller, the pressing roller and the like) of fixing unit 17, as indicated by the two-dot chain lines in FIG. 8, the jammed paper P' is then present between registration sensor S1 for detecting the leading edge of the paper and discharged-sheet sensor S2 for detecting the completion of sheet discharge. The CPU, therefore, cannot detect the presence of the jammed paper P' within the main body of the apparatus when the door switch is closed by closing the cover again, or when the power supply is turned on or off. Thus, the CPU resets the jammed state, and the apparatus moves to the pre-rotation driving in spite of the presence of the jammed paper P'.\nDue to the pre-rotation driving, an accordion-jammed paper remaining in the main body of the apparatus may be more severely jammed, or a jammed wound paper P' may be further wound in the fixing member, making jam removing processing by a user difficult or impossible. When the fixing unit 17 is a film-heating-type unit as in the present case, the thin fixing film 20 may be broken."}
{"text": "The proper care of medical patients is essential for optimal treatment of their medical conditions. Typically, a patient having a particular condition/ailment is prescribed a medicine or treatment based upon established treatment guidelines. The treatment guidelines outline, inter alia, the specific dosages of medicines, the frequency in which dosages should be administered, instructions on how dosages should be administered and the time-lines for therapeutic treatments. Oftentimes, treatment for new patients is administered directly from the treatment guidelines with little variation. These guidelines are typically derived through prospective medical studies. Prospective medical studies, namely, randomized clinical trials, are studies wherein researchers empirically test hypotheses in near ideal conditions by screening the patient population, ensuring that patient care diligently follows the guidelines and recording all relevant data. Such practices fail to take full advantages of historical medical data, rather, relying only on success rates for the patients that rigidly adhered to the treatment guidelines. Additionally, clinical trials are very expensive to conduct.\nHistorical medical data represents a valuable source in the analysis of the patient care process and medical outcomes. As indicated, treatment guidelines have been generated based solely upon the results of treatment on patients who rigidly adhered to the treatment guidelines. However, the number of variables from patient and professional medical care having an impact on the results of patient care is exceedingly high. Moreover, the relationship between these variables is virtually unknown. Accordingly, the ability to fully learn from past medical data could greatly improve patient health care.\nRetrospective studies, for example, the analysis of historical medical patient records from a hospital, are complementary to prospective clinical trials. Health-care organizations are accumulating vast stores of patient data, which are a vital tool for knowledge management. Analyzing this already-collected information may lead to insights that can be subsequently verified in a prospective trial. Most importantly, retrospective studies can measure, in a least two ways, the impact of guidelines in real-life clinical settings. First, retrospective studies can determine the effectiveness of the treatment for a patient population that was excluded from clinical trials. For example, patients above 65, or those with other diseases may be excluded in a clinical trial—however, the guideline validated in that trial is now used to treat all hospital patients. Second, patient treatment in a hospital may differ from that in a trial. For instance, the colon cancer guideline mandates commencing chemotherapy within 6 weeks of surgery, which is rigorously enforced in the clinical trial. However, in a hospital, some patients may begin chemotherapy up to 10 weeks after surgery (e.g., they may be too sick or miss appointments). The impact of this delay on a patient's outcome can only be determined via retrospective analysis since it is not ethical to conduct a clinical trial that would test the impact of this delay—in effect, withholding the accepted standard of care.\nHowever, analyzing hospital data is hard for many reasons. First, medical data is very complex to analyze because of its rich structure. Many traditional statistical methods are ill-suited to data with structure, time-sequenced events (medical data has important temporal components) and/or no structure such as free text, images, etc. Second, because the hospital patient data was collected to treat the patient (as opposed to collected for analysis in a clinical trial), it is imperfect in many ways, for example, missing/incorrect/inconsistent data; key outcomes/variables not recorded; bias in data collection, e.g., sick patients get more tests than well ones, (this is perfectly natural from the medical point of view, but has inherent assumptions that may cause problems for many algorithms); and variables collected/treatments change over time, which particularly impacts some long-term diseases whose treatment can span decades. Lastly, there is wide variation in practice among medical professionals determining if a patient is on a guideline and treated properly is difficult to tell.\nIn view of the above, there exists a need for techniques to collect population-based patient information from a variety of sources, to perform outcome analysis on the collected information, and to conduct retrospective analysis on a large quantity of medical information derived from various sources in a rapid manner."}
{"text": "Technical Field of the Invention\nThe present invention relates generally to accent lighting systems and more particularly to accent lights for use in connection with posts, post caps, stairs and vertical, sloped and horizontal surfaces associated, for example, with decks, docks (and other marine applications), patios, arbors and indoor/outdoor spaces.\nDescription of Related Art\nFor both safety and aesthetic reasons, there is a desire to provide accent lighting in decks, docks (and other marine applications), patios, arbors and indoor/outdoor spaces. For example, stair installations, such as with the treads or risers, are often lit not only for the obvious safety concerns, but also because such lighting is viewed as attractive and opulent. Similarly, the perimeter of an indoor/outdoor space is often accent lit to provide a clear indication of the barrier location but also to provide pleasing illumination of surfaces without having to use overhead lighting which oftentimes is harsh and overly extensive. Still further, accent lighting is often preferred to overhead lighting because the illumination can be better focused and thus the distracting and perhaps visually interfering overspill associated with overhead lighting sources is avoided. Indeed, accent lighting is preferred in many installations, both indoor and outdoor, because such lighting can effectively address safety concerns while simultaneously providing an attractive lighting feature without concern for light pollution.\nA number of accent lighting solutions are known in the art. These solutions suffer from a number of well known drawbacks including excess cost, difficulty or inability for customization, failure of the light source to be hidden from casual view, and premature failure of the lighting sources due to overheating or environmental intrusion (from water, for example).\nThere is a need in the art address the foregoing and other issues when providing an access lighting system."}
{"text": "Network devices, such as routers, are sometimes used to process, route, and provide services to data flows between user devices and server devices. Network devices sometimes include a wide area network (WAN) interface card (WIC). The WIC may include processing instructions, routing tables, or some other information that provide the network device with instructions for processing data flows. Modifying the functionality and/or operations of network devices (e.g., modifying the processing instructions, etc.) may require replacement of WAN cards and/or specialized software development."}
{"text": "1. Field of the Invention\nThe present invention relates to a technique for controlling the screen size of a monitor adapted for the environment of a graphic user interface (hereinafter, referred to as a \"GUI\") between a computer and a user and, more particularly, to a technique for controlling the screen size of a monitor adapted for the GUI environment which is contrived to magnify a part of the screen by using an S-shaped correcting an switching circuit and a correcting port of a horizontal and vertical signal processor in the monitor.\n2. Discussion of Related Art\nIn general human beings can sense graphic data, that is, visual information more rapidly than any other sort of information. For that reason, it may be good if software is visible as other visible objects. A dialogic button controller of Windows utilizing this principle converts conceptual data into visible data. Therefore, data are easily shown when those windows are used, and moreover lots of data can be seen by overlapping the screens. Software utilizing windows is designed to select command words logically by operating the button controller on the monitor screen.\nFor these reasons, as hardware and software have gradually gotten more complicated and their operating method gets more difficult, a GUI (graphical user interface) operational system has been developed. The GUI operational system provides a computer operation with continuity and convenience as well as simplicity of study.\nThose applied programs utilizing the GUI environment can execute more than one program at a time. The applied programs are utilized by dividing CPU time in order. But, this means dividing the time of CPU use, and not simultaneously executing those applied programs. This is because the user cannot sense the speedy conversion of the applied programs by the CPU. The ability to simultaneously perform various functions is called multitasking. An environment for executing the multitasking is called a GUI environment, and is widely distributed recently.\nBut, a monitor adaptable to the GUI environment has not been developed yet.\nIn an earlier monitor for multitasking, the video card receives data processed in a central processing unit of a computer and processes the data into video signals including horizontal and vertical synchronization signals and RGB (red, green and blue) video signals.\nThe RGB signals are fed to a video preamplifier which in turn is fed to a video output amplifier which in turn feeds a cathode ray tube.\nThe horizontal and vertical synchronizing signals are fed to a microcomputer as well as to a horizontal and vertical oscillator signal processor. The horizontal and vertical oscillator signal processor provides an output to a horizontal driving circuit which in turn drives a horizontal output circuit driving a horizontal deflection yoke. The horizontal and vertical oscillator processor also provides an output to a vertical output circuit which drives a vertical deflection yoke. An automatic frequency control circuit provides control signals both to the horizontal and vertical oscillator signal processor and the horizontal output circuit.\nAn S correcting circuit receives an S-shaped correction signal from the microcomputer and supplies the S-shaped correction signals to the horizontal output circuit to correct the distortion of the monitor screen.\nIn such a monitor, during multitasking in the GUI environment, various programs are executed at the same time by using multiple windows. However, the monitor does not provide a circuit which magnifies a certain portion of on the monitor screen using hardware."}
{"text": "This invention relates generally to novel gaskets having continuous multiple seals created by utilizing a core of functionally corrugated material encapsulated by a graphite material such that an interaction occurs among the graphite, the functionally corrugated core, and the surfaces to be sealed. Specifically, the present invention relates to fail-safe, ring-type, concentrically corrugated, graphite encapsulated gaskets with continuous, multiple seals creating a multiplicity of parallel, fluid-locked, graphite-densified, barriers for use in the petrochemical industry to seal flange connections. The present invention also relates to fail-safe, concentrically corrugated, graphite encapsulated heat exchanger and internal combustion engine head gaskets with similarly created continuous, multiple seals.\nIn the petrochemical industry there exists a need for a versatile flange gasket and heat exchanger gasket having the following characteristics: continuous, fail-safe operation thereby avoiding potential damage to human health, plant equipment and the surrounding environment; fire resistance; resistance to temperature differentials across the diameter of the gasket; chemical resistance; field ruggedness; positive sealing even when the gasket is misaligned; positive sealing when the gasket material differs in thermal expansivity with respect to the surface to be sealed; environmentally safe, non-asbestos construction; operability at high pressures and temperatures; resiliency, springiness, memory, and structural integrity particularly during thermal cycling from high temperatures to low temperatures and vice versa, thermal shock, or other thermal excursions; low torque requirements for maintaining a positive seal at all operational or accidental temperatures and pressures to avoid subjecting the flanges (or other surface) and flange bolts to warpage and/or other damage; inhibited gasket creep or cold flow; crash resistance to enhance the potential for gasket reusability; minimized extrusion of the gasket coating; and reduced installation and maintenance time.\nFurthermore, it is likewise desirable to have a cylinder head gasket having similar characteristics as those set forth above in order to improve engine performance.\nUntil the gaskets of the instant application, these desired gasket characteristics, in combination, had not been achieved. Whereas existing gaskets to date have demonstrated a level of utility and selected ones of the characteristics noted, no gasket until the present invention has been able to achieve the desired combination of gasket characteristics.\nGraphite has long been known to possess properties generally desired for high temperature packing materials. In particular, an expanded form of graphite material such as the flexible expanded graphite material sold under the names \"Grafoil.RTM.\" (Union Carbide Corporation, U.S. Pat. No. 3,404,061), \"Calgraph.RTM.\" (Polycarbon, Inc.), \"Flexicarb.RTM.\" (Flexicarb, Inc., a Flexitallic Gaskets, Limited Company), or \"Sigraflex.RTM.\" (Sigri GmbH); an inorganic binder-containing expanded graphite material, such as that described in U.S. Pat. No. 4,234,638; or an intercalated graphite material, such as that described in U.S. Pat. No. 4,799,956, have been employed for use in flange gaskets, heat exchanger gaskets, and automobile engine and exhaust gaskets, and the like.\nBare metal corrugated ring gaskets are known in the art for use as flange gaskets for high temperature applications. Bare metal flat gaskets are also known in the art. However, these bare metal gaskets must be precisely aligned, and the flange and gasket surfaces must be virtually flawless (i.e., highly polished and free of scratches and other abrasions) and the flange bolts must be torqued extremely tightly (approximately 150 ft-lbs on a 3\", 150 lb flange) to create a positive seal. Therefore, the flange bolts are typically tightened and over-tightened with a relatively high level of torque to achieve positive sealing. Such over-tightening can warp the flanges, damage both the flanges and the bolts, and create the dangerous potential for seal failure. The bare metal corrugated gaskets have spring-like characteristics which give them high memory and structural integrity, in contrast with a flat metal gasket which has no memory or resiliency. However, the memory capability of a bare metal corrugated gasket is limited in extreme temperature or thermal cycling situations. Under conditions such as thermal cycling, thermal excursions, or thermal shock, these bare metal corrugated gaskets, and the other existing gaskets tend to leak. Furthermore, bare metal is highly vulnerable to corrosion, abrasion, scratching or other attack from hostile materials thereby creating potential leak paths and gasket failure.\nSpiral-wound gaskets also exist in the art. These gaskets are composed of two materials, for example, a thin profiled metal strip wound into a spiral with an interleaved strip of a sealing material such as Teflon.RTM., paper, asbestos, or graphite. Spiral-wound gaskets have disadvantages. For example, spiral-wound gaskets have inadequate temperature resiliency when subjected to extreme thermal shocking, excursions or temperature cycling. Spiral-wound gaskets are also abusive to the flanges because they are very difficult to seal and require a relatively high level of torque (as much as 150-160 ft-lbs) be applied to the associated flange bolts. Attempts have been made to enhance the elasticity of a spiral-wound gasket, such as by using the gaskets in tandem, or, as suggested in U.S. Pat. No. 4,796,351, by using annealed high tensile strength, precipitation heat-hardenable alloy steel strip and then heat treating such strip prior to interleaving with the sealing material.\nAsbestos gaskets have existed in the art. For example, flat metal, grooved-faced gaskets have been utilized with flanges wherein alternating grooves are filled with asbestos. Asbestos has also been used as a filling material between corrugated sheet metal, metal jackets or spirally wound strips. Also, compressed asbestos sheet material utilizing a rubber binder has been employed in the petrochemical industry as a flange gasket. This asbestos sheet material (such as that sold under the name \"Durabla.RTM.\" by Durabla Manufacturing Co.) is cut into the desired gasket shape and placed between the flange surfaces. Asbestos fiber has desirable material characteristics for gaskets in that it is impervious to most chemicals and acids, it is wear-resistant, and it has high, but limited, temperature capacity (up to approximately 1,000.degree. F. depending upon the stability of the binder material). However, asbestos gaskets only have limited springiness or memory, thereby increasing the likelihood of seal failure during temperature cycling. Furthermore, due to the environmental health dangers associated with asbestos, such a material is no longer desirable in the industry. Kevlar is the most common replacement for asbestos as a gasket material, but it has a temperature limit of only approximately 450.degree. F.\nIt is also known in the art that graphite can be applied to the edge or center of gaskets to improve the sealing characteristics of the gasket. For example, a strip of expanded graphite tape has been applied around the edge of a flat metal gasket, or has been utilized directly on the flange surface. As described below, the use of a graphite tape is undesirable. Graphite has also been used to encapsulate thin, flat metal gasket reinforcements, such as those described in U.S. Pat. Nos. 4,333,975 and 4,422,894. However, these flat metal-reinforced graphite gaskets have no resiliency or memory, lack a plurality of seals, and do not possess a means of preventing seal loss due to extrusion of the graphite from the metal reinforcement during operation at high temperatures and pressures. Furthermore, the metal is often perforated or tanged thereby creating leak paths.\nThe use of a strip or strips of narrow graphite tape (i.e., typically less than one inch width) has been employed on corrugated metal pipe flange gaskets typically having larger than an 18\" diameter. Application of this tape results in covering only a partial planar surface extending from one corrugation apex to an adjacent corrugation apex on only a portion of the corrugated gasket. Thus, the tape surface generally remains flat and parallel to the flange surface, adhering only to the apex portion of the corrugation, but does not extend into the valley areas between each apex. Furthermore, application of this tape to a circular strip of material creates irregularities or flutes in the tape as the tape is bent or folded to conform to the circular shape. This partial application of tape leaves the metal gasket exposed to corrosion, involves non-uniform, non-unitary, distorted, and overlapping application of the graphite to the gasket face, and does not create a plurality of concentrically spaced, parallel, fluid-locked barrier seals. Furthermore, when the adhesive backing carbonizes at high temperature, and the gasket is placed through a series of temperature cycles, the resulting expansion and contraction of the metal can thereby cause stress and loss of adherence at the graphite/apex points of contact, and/or cause the graphite layer expanse between adjacent corrugation apex points to tear. Furthermore, the expanse of graphite material between adjacent corrugation apex points is vulnerable to puncturing, tearing, and cutting during handling, installation, and use. Thus, the application of this tape creates leak paths, and the potential for extrusion of the graphite material and seal failure.\nIt has also been suggested that a thin rectilinearly embossed iron sheet could serve as a reinforcement for an inorganically bound graphite material, such as in U.S. Pat. No. 4,234,638. However, this thin rectilinear embossment serves only as a thin reinforcement to a thicker graphite seal, and is not believed to be interactive with the surface to be sealed, nor does it prevent extrusion or channelling of the graphite. This arrangement may experience seal failure at high temperature and pressure.\nGraphite has also been used with grooved flat metal pipe flange gaskets, such as described in U.S. Pat. No. 4,121,858. However, the grooves in the gaskets function only to help contain the graphite, and are not believed to give the gasket any memory or resilience.\nFurthermore, graphite has been used on heat exchanger head gaskets (for example, as a filler material on a double jacketed heat exchanger gasket as set forth in U.S. Pat. No. 4,118,850 and as a single flat surface coating on a heat exchanger gasket as described in U.S. Pat. No. 4,872,506). However, these heat exchanger head gasket seals also require that a relatively high level of torque (up to approximately 150 ft-lbs) be applied to the head bolts to achieve a positive seal, and these heat exchanger gasket seals have no structural memory or resiliency.\nGraphite has also been employed in conjunction with the manufacture of internal combustion engine head gaskets and exhaust manifold gaskets, such as those described in U.S. Pat. Nos. 3,926,539; 4,083,570; 4,423,544; 4,465,287; 4,591,170; 4,705,278; 4,723,783; 4,776,602; 4,810,454; 4,822,062; and 4,911,972. Automotive head gaskets typically require that about 90-100 ft-lbs of torque be applied to the associated bolts to achieve a positive seal and are essentially constructed from a flat, thin (e.g., 7 mil.) metal sheet material which has been perforated and subsequently laminated with a graphite material. The metal has perforations and may further include protrusions (or tangs) extending outwardly on one or both sides of the face of the metal. The graphite material is applied to the surface of this perforated or tanged metal and then the gasket is compressed to close the tangs to affix the graphite in place. Problems exist with these gaskets. For example, the perforations and tangs in the metal substrate may provide leak paths. Additionally, the actual perforated metal is believed to be too thin for welding, thereby, creating size limitations for the gasket. Even if one could weld this very thin (e.g., 7 mil.) material to create a large diameter gasket, the material must then be tanged, which again presents difficulties due to the increased size. Furthermore, these tanged, flat gaskets do not have resilience or memory, thereby reducing their effectiveness in high performance, temperature cycling applications.\nGraphite has also been employed in conjunction with the manufacture of valve stem seals or stuffing box seals, such as those described in U.S. Pat. Nos. 4,116,451; 4,160,551; 4,190,257; 4,350,346; 4,394,023; 4,744,572; 4,753,443; and 4,892,320.\nAs noted above, these existing gasket configurations have limitations and are inadequate to accomplish the desired, combined characteristics for flange, heat exchanger and engine head gaskets as set forth herein.\nTherefore, in accordance with the present invention, a fail-safe, multi-sealed, interactive gasket system is provided which overcomes or substantially minimizes these and other disadvantages of the gaskets presently known or used in the art."}
{"text": "In a service chain scenario, after a service path is specified for some traffic, the traffic passes through a flow distribution node and each service node according to the service path and a specified sequence. However, when packet processing is performed between adjacent nodes of the flow distribution node and service nodes, flow rules between the adjacent nodes may be the same, or an upstream node can support a flow rule of its downstream node, and packet processing performed by the downstream node such as discarding, allowing passing, and forwarding, can also be performed by the upstream node. In many cases, and especially in a case in which flow processing by the upstream node is implemented using hardware, how many actions are performed on a flow is irrelevant to traffic processing performance of the upstream node, which is flow processing performance in brief, or how many actions are performed does not have great impact on flow processing performance. That is, no matter how many actions are performed on the flow by the upstream node, the flow processing performance of the upstream node does not change greatly, for example, flow processing implemented using an application-specific integrated circuit (ASIC). In the foregoing cases, if flow processing performed by the downstream node is deployed in its upstream node, burden of the downstream node can be reduced, and there is no great impact on the flow processing performance of the upstream node. In the foregoing content, the upstream node and the downstream node are two adjacent nodes.\nGenerally, flow processing by the downstream node may be deployed in its upstream node by performing traffic offloading on the downstream node.\nHowever, in the prior art, no specific solution to traffic offloading on the downstream node is provided. Therefore, it is very difficult to perform traffic offloading on the downstream node."}
{"text": "1. Technical Field\nThe present disclosure relates to a travel control device for a vehicle including a steering transmission ratio variable device and a power steering device.\n2. Description of the Related Art\nAs a travel control device for vehicles such as a motor vehicle, there has hitherto been well known a travel control device including a steering transmission ratio variable device and an electric power steering device in a steering transmission system between a steering wheel and steered wheels. The steering transmission ratio variable device is arranged between an upper steering shaft and a lower steering shaft, and is configured to rotate the lower steering shaft with respect to the upper steering shaft, thereby changing a steering transmission ratio. This type of travel control device is allowed to control a steering angle of the steered wheels without requiring a steering operation by the driver, thereby being capable of optimally controlling the steering transmission ratio, namely, a ratio of a change amount in a steering angle of the steered wheels with respect to an operation amount of the steering wheel depending on a vehicle speed or the like.\nFurther, there is also known a travel control device including a driving support device configured to calculate a target steering angle of the steered wheels for causing the vehicle to travel along a predetermined travel path, and to control the electric power steering device and the steering transmission ratio variable device so that an actual steering angle of the steered wheels conforms to the target steering angle. The travel control device including the driving support device may cause the vehicle to travel along the predetermined travel path, and may attain desired steering characteristics without requiring the steering operation by the driver.\nFor example, in Japanese Patent Application Laid-open No. 2011-31770, there is described a travel control device configured to carry out driving support control of controlling the electric power steering device, to thereby control the steering angle of the steered wheels so that the vehicle travels along a predetermined travel path, and to suppress a behavior change of the vehicle generated during the driving support control. With the travel control device described in Japanese Patent Application Laid-open No. 2011-31770, the vehicle can be caused to travel along the predetermined travel path while suppressing the behavior change of the vehicle.\nHowever, there are the following problems in a related-art travel control device in which the electric power steering device is installed in the steering transmission system between the steered wheels and the steering transmission ratio variable device like the travel control device described in Japanese Patent Application Laid-open No. 2011-31770. Specifically, when the electric power steering device is controlled through the driving support control and the steering angle of the steered wheels is changed, a steering member like a lower steering shaft between the electric power steering device and the steering transmission ratio variable device is rotated. A friction force acts between members rotating relatively to each other in the steering transmission ratio variable device. Thus, when those members cannot rotate relatively to each other due to the friction force, an upper steering shaft on a steering wheel side and the lower steering shaft on a steered wheel side with respect to the steering transmission ratio variable device do not rotate relatively to each other, but integrally rotate.\nA steering angle sensor configured to detect a steering angle required to control the steering transmission ratio and the like is installed on the upper steering shaft, and thus the steering angle detected by the steering angle sensor changes as the upper steering shaft rotates. The steering transmission ratio variable device is controlled so as to attain a predetermined steering transmission ratio based on a changing steering angle, and thus the steering transmission ratio variable device tries to rotate the lower steering shaft with respect to the upper steering shaft even when the driver is not steering. However, the electric power steering device is generating a steering torque so as to change the steering angle of the steered wheels, and thus the steering transmission ratio variable device cannot rotate the lower steering shaft, while the upper steering shaft is rotated by a reaction force.\nIn particular, while the driving support control is being carried out, the driver does not need to actively carry out the steering, a force of the driver for gripping the steering wheel is low, and a so-called hands-off state may be brought about. Therefore, a repetition of the above-mentioned phenomenon causes a repetition of a reciprocal rotation of the upper steering shaft. As a result, a minute vibration of the steering wheel and a vibration change in a steering reaction force are liable to be generated, and the driver may feel a sense of discomfort and a sense of annoyance.\nIn the travel control device described in Japanese Patent Application Laid-open No. 2011-31770, as shown in FIG. 8 of Japanese Patent Application Laid-open No. 2011-31770, when the driving support control is being carried out, the steering transmission ratio is changed toward a small side (so-called slow side) compared to a case in which the driving support control is not being carried out over an entire vehicle speed range. Thus, influence of the steering operation by the driver exerted on the change in the steering angle of the steered wheels can be reduced, thereby increasing a travel stability of the vehicle, but the problem of the minute vibration of the steering wheel and the vibration change in the steering reaction force cannot be solved."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for holding springs of a clutch, particularly to shock-absorbing springs for a clutch plate which are used in a direct-coupling clutch of a fluid torque converter, or the like.\n2. Description of the Related Art\nA direct-coupling clutch for a fluid torque converter is disclosed, for example, in Published Examined Patent Application No. 9780/1987. An apparatus for holding shock-absorbing springs in this direct-coupling clutch has a construction shown in FIG. 12A. In the figure, reference mark \"a\" is a clutch plate. On a front side (i.e., left side in FIG. 12A) of the clutch plate \"a\", there is adhered a friction member b. On a rear side (i.e., right side in FIG. 12A) thereof, there is formed an annular recessed portion c. A shock-absorbing spring d is held inside this recessed portion c by means of a wear-resistant guide plate e and a retainer f. The retainer f has end surface holding portions f.sub.1 for holding both end portions of the spring d, and periphery holding portions f.sub.2 for holding the periphery thereof. The retainer f is fixed to the clutch plate \"a\" with rivets g after overlapping it with the guide plate e. In the figure, reference mark h is an input case, i is a turbine, and i.sub.1 are projection pieces which extend from the turbine i to abut with both ends of the spring d.\nIn this apparatus, as shown in FIG. 12B, prior to fixing with rivets g, the guide plate e and the retainer f are assembled in advance by caulking a small projection f.sub.3 which protrudes or projects from the retainer f into a small hole e.sub.1 which is provided in the guide plate e. However, the guide plate e is likely to give rise to distortions because it is made of a thin plate and is further heat-treated. Therefore, the portion which is overlapped with the retainer f presents remarkable distortions because it is planar and is wide in lateral width. Consequently, it sometimes happens that a clearance occurs between the guide plate e and the retainer f in a portion away from the portion that is caulked by the small projection f.sub.3. Even if the pre-assembled parts are caulked onto the clutch plate \"a\" with rivets g as shown in FIG. 12A, the above-described distortion will remain, resulting in a disadvantage in that there will be left a clearance between the guide plate \"a\" and the retainer f or the clutch plate e.\nFurther, in order to minimize the number of parts, it is considered to make the arrangement, as shown in FIG. 13, that a small projection j is provided in a clutch plate \"a\" and that the guide plate e and the retainer f are caulked at the same time without pre-assembling them. However, it is considerably difficult to protrude or project the small projection j into a height enough to caulk the guide plate e and the retainer f together. In addition, the guide plate e may be distorted to thereby increase in height. As a result, the height of the small projection j is likely to become insufficient in plate thickness at a caulking head portion j.sub.1. Further, due to the distortion in the guide plate e, the clearance as described above is likely to remain between respective plates.\nThe present invention has an object of providing an apparatus for holding a spring of a clutch in which the retainer can be mounted on the clutch plate without giving rise to a clearance."}
{"text": "It is often necessary for certain telephone-specific or switching system-specific operating steps to be performed in order to program a function key. The allocation of a programmed speed dialing key to a telephone number is usually indicated here by the user manually writing on a paper label strip of the telephone with a pen.\nThere are also modern telephones whose keys are assigned an LCD or LED display (LED: Light-Emitting Diode; LCD: Liquid Crystal Display) in each case, on which the stored telephone number of the programmed speed dialing key can be displayed as an electronic key label. Manual labeling of the keys is not required for such telephones. Function keys of this type with a display are often referred to as self-labeling keys.\nIn known communication systems, setting the display for a key is then carried out during programming of the function key. For instance, following programming of a speed dialing key, the telephone number allocated to the speed dialing key is displayed on the display for the speed dialing key. In addition, as part of programming the speed dialing keys, the user is often able to manually edit the electronic key label via a user interface, by entering a text or the name of the subscriber for the stored destination telephone number for example. This name can then be displayed on the display of the speed dialing key instead of the telephone number."}
{"text": "In the area of conventional business service development, there are, in general, three major roles, each role having specific tools to address its needs.\nApplication architects are primarily concerned with defining the contracts of the business services to meet the business requirements. An application contract consists of one or more interface definitions, reflecting the right level of interface granularity and service behavior such as operations exposed and their inputs and outputs.\nAn important aspect of the architect's task here is to ensure that existing interface types are reused so that multiple interface types and messages are not reinvented to fulfill similar requirements.\nOnce the correct interfaces have been defined, the application developers are ready to implement those interfaces. This often involves writing code that implements business logic. The developers also have to be able to unit test the code that they have just written.\nFinally, once the code has been developed, it can be given to the system administrators for configuration, deployment and management. They can then configure the runtime execution environment for the newly developed service and deploy the application. Once deployed, the application can then be managed through the means available in the environment.\nIn the conventional development model, the applications are developed from scratch. This allows the users to perform top-down design and leverage the full capabilities of the available technologies since no baggage needs to be carried from the past.\nCurrently, there are 5 major phases of the lifecycle of the business services. In the model phase, the application architects design the interfaces for the business services to be developed. Once the interfaces have been designed, the business logic is developed by the developers of the system. The system architects then configure the runtime for the business service and the administrators then deploy and manage the business services.\nIn many cases, the business services already exist but are not manageable. These services may be either 3rd party off-the-shelf (COTS) applications or custom applications developed in-house or by system integrators.\nThis style of development is bottom-up development where the interfaces are actually constructed based on existing applications. The architects first use introspection to reverse engineer metadata from the existing applications. These applications may exist in one of many supported technologies such as Java classes, CORBA, EJB's, existing web services or COTS packages such as SAP or Siebel. Introspection would result in one-to-one interface generation. However, in real world, these interfaces may need to be refactored. Refactoring involves either suppressing some of the operation in an interface, creating brand new ones, combining two or more interfaces (contract aggregation) or splitting a contract into two or more contracts (contract dissemination).\nIt should be noted that the existing applications could also be web services. The process should however be no different than that for 3rd party COTS or custom applications supporting a particular technology.\nBusiness logic generally does not change in nature. However, the underlying technologies are changing very rapidly. For example, there are a number of new standards that are being developed to address various aspects of distributed business services. Many standards do not have any sort of industry consensus and many areas have competing standards. However, the applications have to be developed today. Typically, the business services are written to provide all the other necessary functionality, such as security and others (application infrastructure services) that do not actually contribute to the business behavior itself. As the standards and technologies for these other services (application infrastructure services) evolve, the business services have to be modified and enhanced in order to take advantage of them. Furthermore, in the current state of the art, the configuration and deployment of these services is specific to a particular underlying environment and the underlying products.\nIn this type of distributed system, it is important that the deployed packages are fully supported by the underlying processes, and that the packages can be easily configured and deployed for use. Often, specific business services require multiple infrastructure services, and it is a common, costly, and time-consuming process to customize business service packages with all the required underlying services.\nIt would be desirable to provide a system, method, and computer program product which would allow a system to automatically generate program code that insures proper interaction of business services and infrastructure services."}
{"text": "Emulsion aggregation (EA) toners are used in forming print and/or xerographic images. Emulsion aggregation techniques typically involve the formation of an emulsion latex of resin particles that have a small size of from, for example, about 5 to about 500 nanometers in diameter, by heating the resin, optionally with solvent if needed, in water, or by making a latex in water using an emulsion polymerization. A colorant dispersion, for example of a pigment dispersed in water, optionally with additional resin, is separately formed. The colorant dispersion is added to the emulsion latex mixture, and an aggregating agent or complexing agent is then added and/or aggregation is otherwise initiated to form aggregated toner particles. The aggregated toner particles are heated to enable coalescence/fusing, thereby achieving aggregated, fused toner particles. United States patent documents describing emulsion aggregation toners include, for example, U.S. Pat. Nos. 5,278,020; 5,290,654; 5,308,734; 5,344,738; 5,346,797; 5,348,832; 5,364,729; 5,366,841; 5,370,963; 5,403,693; 5,405,728; 5,418,108; 5,496,676; 5,501,935; 5,527,658; 5,585,215; 5,650,255; 5,650,256; 5,723,253; 5,744,520; 5,747,215; 5,763,133; 5,766,818; 5,804,349; 5,827,633; 5,840,462; 5,853,943; 5,853,944; 5,863,698; 5,869,215; 5,902,710; 5,910,387; 5,916,725; 5,919,595; 5,925,488; 5,977,210; 6,576,389; 6,617,092; 6,627,373; 6,638,677; 6,656,657; 6,656,658; 6,664,017; 6,673,505; 6,730,450; 6,743,559; 6,756,176; 6,780,500; 6,830,860; and 7,029,817; and U.S Patent Application Publication No. 2008/0107989.\nThe disclosures of each of the foregoing patents and publications are hereby incorporated by reference herein in their entireties. The appropriate components and process aspects of each of the foregoing patents and publications may also be selected for the present compositions and processes in embodiments thereof.\nCurrent formulations of toner show a need for improved fusing performance. Poor fusing creates problems in paper adhesion and print performance.\nCurrent formulations of toner show a need for improved flow. Poor flow creates problems in gravity-fed cartridges, causing toner to hang up due to poor flow properties and leads to deletions on paper.\nThere remains a need for an improved toner composition and process that overcomes or alleviates the above-described and other problems. There further remains a need for a toner composition suitable for high speed printing, particularly high speed monochrome printing that can provide excellent flow, charging, lower toner usage, and reduced drum contamination."}
{"text": "Firearms use rapidly expanding gasses, typically from an explosive charge, such as smokeless powder, to accelerate a projectile down a barrel of the firearm toward a target. In contrast to bullets used in many rifles and handguns, shotgun shells include a collection of “shot” in the shell. The shotgun uses an explosive charge to accelerate a shot toward a target. Shot can take various forms, including various sizes, quantities, packing orientations, shapes (e.g., spherical, cubic, tetrahedron, etc.), and compositions. The shot may be initially contained by a wad during acceleration of the shot in the barrel. The wad exits the barrel of the shotgun with the shot while the shell remains in the shotgun. Conventional wads may travel a short distance with the shot, but are designed to experience atmospheric drag such that the shot separates from the wad in flight.\nShotguns utilize shot instead of a bullet to allow the shot to spread over an area as the shot travels from the barrel to the target. A bullet is intended to remain a single object during flight of the projectile to the target. The spread of the shot at the target may depend, at least partially, upon a number of characteristics of the shotgun and the shell fired, such as the muzzle velocity of the shot, the length of the barrel, the type of shot, and the distance to the target. The shape and dimensions of the area over which the shot spreads during the flight of the shot is known as the “pattern” of the shot. The pattern may be important to a shooter, as different patterns may be desirable for different purposes and different types of shot.\nThe shotgun itself may be altered or customized to modify the pattern of the shot. For example, the barrel may be shortened and/or widened (e.g., a home defense shotgun) to decrease the density of the shot pattern (i.e., create a larger spread to the pattern) at the expense of effective range and velocity of the shot. In contrast, a barrel may lengthened and/or constrained (e.g., a choke may be added) to increase the velocity of the shot and to increase the density of the shot pattern (i.e., reduce the area over which the shot spreads).\nThe wad may also affect the pattern. The wad will exit with the shot after moving through the barrel. The wad is lighter than the shot and will decelerate from drag with the air more readily than the shot, causing the shot and wad to decouple during flight to the target. Rotation of the wad during decoupling may cause the pattern to deviate in the direction of the rotation. Conventional wads are designed to decouple quickly from the shot to minimize the impact unintended rotation may have on the direction and, hence, pattern of the shot. However, after the shot is released from the wad, the shot may experience additional turbulence and drag in the air, causing the shot to slow and spread, reducing the effectiveness of the shot at distance."}
{"text": "This invention is related to the composition of an elastomeric triblock copolymer, and, additionally, to the use of such an elastomeric copolymer to prepare dip-formed articles.\nThin-walled, extensible articles such as gloves and condoms have long been made from natural rubber. In normal production, such articles are formed from natural rubber latex, a naturally occurring emulsion of rubber and water, with added stabilizing agents and vulcanizing chemicals. A form of the appropriate shape, previously coated with a coagulating solution in some cases, is dipped into the latex mixture once or several times to build up a layer of the desired thickness. The water is allowed to evaporate, leaving a solid rubber film. The film must be vulcanized to provide adequate mechanical and physical properties.\nNatural rubber has many advantages in these applications, being strong and highly elastic and having good \"tactility\" or feeling to the user. The good tactility is believed to be a result of its low deformation stress at 10-500 percent elongations and high elastic recovery from these elongations. Natural rubber has the shortcoming that the resulting rubber film devices are sometimes prone to having \"pinholes\" therethrough, often on the order of micrometers in diameter. The pinholes are probably due to minute impurities in the natural latex which are difficult to filter out, and to the fact that the process converts a stable emulsion to a film by coagulating chemicals or heat. Natural rubber is also susceptible to rapid attack by ozone which causes scission cracking, as well as to oxidative attack during storage which causes cracking and destroys the physical integrity of the product. Even with these problems, for many applications the pinholes do not prevent the rubber article from having utility, and latex gloves and condoms are widely used. However, the presence of the pinholes is unacceptable in other applications, as for example where there is any possibility of the presence of the AIDS virus. Natural rubber is also not hypoallergenic due to the residual surfactants, vulcanizing agents, stabilizing agents, antioxidants, or protein materials in the rubber. Persons who are particularly susceptible to irritation or sensitization, or use the rubber products for extended periods of time may experience allergic reactions.\nVarious types of synthetic elastomeric polymer products have been developed for use in thin articles produced by dip forming. Synthetic rubber compositions can be dissolved in solvents to form a true solution, so that pinholes are much less likely to be present. Available synthetic rubber compositions have various other shortcomings, including unacceptable tactility. While each may meet some of the requirements, the known synthetic compositions do not have the required combination of strength, tactility, resistance to environmental damage, and hypoallergenicity required for many products such as examination and surgical gloves and condoms. Several synthetic elastomeric systems have provided strength and tactility, but none has solved the problems related to environmental damage and hypoallergenicity required for surgical gloves and condoms.\nThere is therefore a need for an improved elastomeric material for use in thin, dip-formed articles such as examination or surgical gloves and condoms. Such a material must have the required properties of strength and elastic elongation. It also must be pinhole free when the article is formed and used, resistant to immediate environmental or damage occurring during storage or use, and hypoallergenic. The present invention fulfills this need, and further provides related advantages."}
{"text": "The present invention relates to a method of monitoring the movements of an individual in and around buildings, rooms and other spaces, and then particularly a method of enabling the movements of individuals to be monitored during military or civilian exercises."}
{"text": "The determination of the optical properties of tissues is of fundamental importance in many fields of medicine, both for diagnostic and monitoring purposes. It is well known that light of differing wavelengths penetrates differently in various tissues. In the near infrared region, for example, (about 650 nm to about 1000 nm), light of this wavelength penetrates several centimeters through tissue. It is intended that the term \"light\" includes other electromagnetic radiation as well which is invisible to the human eye, for example, infrared and ultraviolet.\nBecause of the capability of various forms of light to penetrate tissue for several centimeters, photometric or spectroscopic methods can be used to measure the concentration of tissue metabolites such as hemoglobin by the measuring of the absorption of the light at one or more wavelengths. It is desirable and important that apparatus which measures the optical properties of tissues for clinical purposes provide quantitative information of a desired parameter, for example the concentration of oxyhemoglobin, or deoxyhemoglobin, glucose, or other metabolites.\nIn normal practice, the absolute determination of the concentration of a substance can be obtained by the measurement of the light transmitted through a sample of known thickness. Such a transmission measurement enables one to determine the absorption coefficient. Using this, the concentration of the measured substance can be calculated using the molar extinction coefficient of that substance via the Beer-Lambert law.\nIn the event of interference caused by more than one substance being present, measurement at different wavelengths can provide a method to determine the concentration of one or more different chemical species present, assuming that the materials present have different absorption spectra. The success of this method depends on the precision of the measurement and on the number of different substances present.\nAdditional problems arise in the photometry of tissues and other materials having high turbidity, such as emulsions. For purposes of this disclosure, it is to be understood that the term \"tissue\" includes living materials, but can also include non-living materials such as emulsions when it is desired to obtain similar data from such emulsions as is done by this invention with tissues. A measurement of the light transmitted through a slab of tissue has in the prior art been not practical, using non-invasive methods, except for special, thin regions of the body where light can shine entirely through the tissue and be detected on the other side. An example of this is a clinically used photometric blood oxygen sensor, which fits on the finger tip and shines directly therethrough to give real time oxygen concentration data.\nIn tissue photometry, the amount of transmitted light depends not only on the absorption of the medium being analyzed, but also on the scattering properties thereof. This light scattering greatly increases the complexity of photometric analysis of tissue, emulsions, and similar materials, since light scattering produces an unpredictable variation of the amount of light transmitted, which can vary significantly between various samples of tissues and the like.\nMany different methods have been proposed to deal with this problem of scattering in photometric processes. For example, empirical corrections based on the type of tissue to be measured have been used to account for the effect of scattering on the absorption properties. For reflection measurements, theoretical models have been used to calculate the albedo of a surface. The success of all of these models has been poor, although there are commercially available instruments based upon those principles. A major problem is that in order to obtain a reasonable estimate of the concentration of a substance in tissue, some sort of a priori calibration must be performed, based on a statistical analysis of a large number of corresponding tissue samples. However, the range of variation of scattering within tissues from various individuals results in fundamentally unpredictable results, with the photometric results being strongly modified by factors such as skin color, and the amount of lipids in muscles.\nThe Hamamatsu Company of Japan in 1990 introduced a simple tissue spectrometer called NIRO 500 for the measurement of tissue oxygenation and total blood volume for neonatal monitoring. The principles of this device are as disclosed in Cope U.S. Pat. No. 5,032,024. The instrument is a steady state instrument, and is based on four different laser diodes emitting in the near infrared range. The light is brought to the tissue using a fiber optic system. The measurement is purely a steady-state one, with the optical path length in the tissue being not measurable. Thus, contrary to this invention, only relative quantities can be obtained, rather than absolute concentrations of oxy- and deoxyhemoglobin.\nBy this invention, absolute quantities of materials found admixed in highly turbid media may have their concentrations determined in a quantitative manner. Specifically, by this invention tissue metabolites may be quantitatively determined in real time, on a continuous basis, for example, concentrations of oxy- and deoxyhemoglobin, glucose, or the like. This can be accomplished without the need to pass light through a narrow portion of tissue, for example an extremity such as the finger. Rather, a sensor may be placed on a more central area of the body for determination of metabolite concentrations or other parameters there. This may be accomplished in a non-invasive manner, essentially instantaneously. Also, it may be possible for different metabolites present to be selectively and quantitatively determined on an essentially instantaneous basis."}
{"text": "Many different types of toasters have been manufactured and used all over the world for decades, ranging from small single or double slice home toasters to large commercial toasters. Toasters may be used to toast sliced bread, rolls, bagels, biscuits, and other similar food products.\nOne type of toaster is a conveyor toaster that continuously moves food prod products through the appliance. Generally, these conveyor toasters include a housing having a toasting cavity, a conveyor for moving the bread/toast through the toasting cavity, and one or more heating elements.\nVarious efforts have been made to develop conveyor toasters, and the following references show a few examples. These references are incorporated by reference herein, in their entireties:\nU.S. Pat. No. 1,696,613: Toaster\nA toaster includes a horizontal conveyor, the bread to be toasted being-placed on the conveyor at the front of the toaster and carried toward the rear in close proximity to suitable heating means, whereby the bread is toasted. At the rear of the toaster the bread is fed by the conveyor to means for returning it to the front of the toaster, this means being shown as an inclined slide. The toaster further includes means for regulating a draft of cooling air around the heating means whereby the generation of heat may be-regulated for different toasting operations and for light or heavy loads, and also a special crumb tray arrangement below the heating means.\nU.S. Pat. No. 1,708,522: Toaster\nA toaster is provided in which the articles to be toasted, such as slices of bread, are fed through the apparatus by suitable conveying means, and wherein the heating of the heating elements is started by insertion of the material to be toasted. The toaster includes conveying means for the material to be toasted and means for driving or stopping said conveyer in accordance with the presence or absence of the toaster of material to be toasted, and further includes a casing having insertion openings at different positions along the conveyer so that the length of path traveled by the bread in passing the heating means may be varied to vary.\nU.S. Pat. No. 3,517,605: Bun Toaster\nA compact and easily cleaned bun toaster such as for the heels and crowns of hamburger buns and the like. A conveyor confronts a toasting platen for driving buns along the surface of the conveyor to toast them. The conveyor is removable as a unit from the frame mounting the platen, thereby to expose both the conveyor and the platen for easy cleaning. The platen is movable toward and away from the conveyor to adjust the space therebetween. The platen is stepped to facilitate simultaneous toasting of bun heels and crowns of different thicknesses. The conveyor is in driven engagement only when moved into suspended engagement with the frame.\nU.S. Pat. No. 3,693,536: Apparatus and Method of Toasting Bread Like Articles\nAn electric toasting apparatus comprised of a conveyor adapted to support articles to be toasted and having the transverse width and a preselected rate of movement in a longitudinal direction. The apparatus includes at least two radiant heater means disposed adjacent the conveyor, each arranged to provide generally uniform heating conditions across substantially the entire transverse width of the conveyor and along at least a substantial portion of the longitudinal length of the conveyor. One heater runs constantly and has a heat output to bring the articles close to the toasting point. The other is modulated to control its heat output and thus the total output of both heater means to give the desired degree of toasting. The method comprises the steps of: (a) conveying the articles past and adjacent to the first and second heater means at a uniform rate; (b) fully energizing the first heater means in an amount sufficient to bring the articles at least close to the toasting point; and (c) simultaneously energizing the second heater means to produce less than its maximum heat output, but sufficiently to fully toast the articles.\nU.S. Pat. No. 3,835,760: Apparatus for Toasting Bread-Like Articles\nApparatus for toasting bread-like articles is disclosed comprising a housing providing a toasting chamber in which electrical heating elements are disposed together with an endless conveyor operable to receive and carry articles to be toasted past the heating elements. The housing has an entrance opening for introducing articles onto the conveyor and when the articles have been toasted they are discharged from the conveyor for removal from the apparatus. The housing is provided with a discharge opening in the top wall thereof and adjacent the rear wall of the housing, and a flue arrangement is provided within the housing to induce hot air in the housing to flow through the discharge opening. Flow of hot air through the flue induces ambient air to flow into the toasting chamber through the product entrance, beneath the conveyor and heating elements and thence upwardly toward the top of the chamber. A temperature sensing element is disposed in the housing in the path of ambient air flow and is responsive to the temperature sensed to cause energization and de-energization of the heating elements.\nU.S. Pat. No. 4,286,509: Toasting Apparatus\nAn economical energy saving toasting apparatus includes a housing with an interior baffle defining a toasting chamber therein and an inlet for the introduction of the product to be toasted into the chamber. The housing also includes an outlet through which the toasted product is dispensed. A conveying assembly for conveying the product from the inlet through the toasting chamber to the outlet is also included within the housing. Mounted within the toasting chamber is a first set of resistive heating elements that are electrically connected in a circuit and continuously energized during the operation of the heating apparatus. A second set of heating elements is also included in the heating chamber and connected to a timer in order to be energized for only selected periods of time. A micro-switch is also included in the apparatus to actuate the timer and a micro-switch actuator is positioned in the inlet and adjacent to the conveying device so as to be engaged by the product prior to entry into the toasting chamber. The engagement of the product actuates the micro-switch to energize the timer and the second set of heating coils for a predetermined period of time.\nU.S. Pat. No. 5,673,610: Apparatus for Conveyorized Toasting of Breads and Like Food Items\nAn apparatus for conveyorized toasting of sandwich buns and like bread and food items simultaneously on opposite sides comprises a central heated platen with two bun transport conveyors traveling in spaced relation along opposite sides of the platen and with a pair of auxiliary heating elements disposed outwardly of the respective food transport runs of the conveyors in facing relation to the opposite sides of the platen. Each conveyor is biased toward the platen by a pair of parallel pivot arms urged by springs into pivoted engagement against the transport run of the respective conveyor to define a predetermined desirable spacing to the facing side of the platen.\nU.S. Pat. No. 5,746,116: Rapid Toasting Apparatus\nAn improved toasting apparatus is disclosed including a variable speed conveyor belt for conveying products through the toaster; a bottom heating element positioned below the conveyor belt, the bottom heating element including variable control for controlling the amount of heat provided by the bottom heating element; a drying zone including a first plurality of upper heating elements positioned above the conveyor; a toasting zone including a second plurality of upper heating elements; and wherein the bottom heating element and the first and second plurality of upper heating elements are controlled individually and independently.\nU.S. Pat. No. 5,983,785: Contact Toaster with Infinite Adjustment\nA contact toaster has a housing with a product inlet, product outlet and at least one toasting chamber defined in the housing. A heated platen is mounted in the housing and has a platen surface arranged for toasting the food products. A flexible endless belt, which when rotated, is arranged to pass the food products in a toasting path along the platen between the belt and the platen. An infinite controller acts to selectively displace the belt to adjust the distance between the belt and the platen to any value between minimum and maximum distance limits to accommodate food products of different widths. A guide is situated between the belt and the infinite controller so that the controller is in contact with at least one area of the guide. The infinite controller includes a variable distance element with a cam surface which is settable to any selected point of the cam surface between minimum and maximum cam radial limits to move the guide and belt to any belt to a selected platen distance that can have any value equal to or intermediate the maximum and minimum distance limits.\nU.S. Pat. No. 6,019,030: Conveyor-Type Food Product Toaster\nA toaster for bread-type food products includes a housing defining at least one toasting cavity. A toasting heater is operatively associated with the cavity. An endless chain-type conveyor is provided for conveying food products through the cavity. An endless belt is disposed about the conveyor to prevent the chain-type conveyor from imprinting the food products. The belt includes a seam, and a retention portion of the belt at the seam is engaged with the conveyor to prevent slippage between the belt and the conveyor.\nUS Publication Number 20100139497A1: Food Heating Device\nA food heating device includes a first conveyor having a first thermally-conductive belt and a second conveyor having a second thermally-conductive belt. A first platen is disposed on a first side of the first thermally-conductive belt and a second platen disposed on a first side of the second thermally-conductive belt. The first conveyor and the second conveyor are arranged and spaced to transport a food product between a second side of the first thermally-conductive belt and a second side of the second thermally-conductive belt.\nUS Publication Number 201000275789A1: Toaster with Removable and Adjustable Conveyors\nA dual-sided, conveyor toaster provides operator-adjustable conveyors that are also operator removable. The removable conveyor assemblies are chain driven and removably supported in the toaster by re-positionable mounting mechanisms embodied as either adjustable pins that engage fixed slots or adjustable slots that engage fixed pins. Fixed pins can be located on the conveyor and engage adjustable slots in the toaster; adjustable pins can be located on the toaster and engage fixed slots on the conveyor. The conveyor assemblies use non-marring plates to urge food products against the heated platen surfaces and to carry the food products across the platen."}
{"text": "The present invention relates to the art of plasma arc torches and, more particularly, to improvements in connection with a drag cup for use with a plasma arc torch.\nIt is of course well known that a plasma arc torch comprises a nozzle receiving an electrode which has a nose end facing an end wall of the nozzle and which end wall has a plasma outlet opening therethrough. The electrode and nozzle may be relatively displaced between starting and operating positions in which the electrode respectively contacts the end wall and is spaced an operating distance therefrom. The end wall of the nozzle and the end face of the electrode provide a gas chamber into which a plasma gas is supplied and from which a plasma jet is emitted through the outlet orifice when an arc current is flowed between the electrode and nozzle. Upon starting the torch, by moving the electrode out of contact with the end wall to create a pilot arc, the torch operates in a non-transferred pilot arc mode and, when the nozzle is moved into proximity with a workpiece, the arc is transferred to the workpiece and the torch then operates in an arc-transferred mode. In an alternative structural arrangement, the electrode and nozzle can be fixed relative to one another and the torch started by the use of a high frequency or other known starting procedure.\nDrag cups are used with hand held plasma arc torches and advantageously allow an operator to contact the torch assembly with a workpiece such as during a cutting operation, primarily to assist the operator in maintaining a correct standoff distance between the torch and workpiece. A drag cup also functions to protect against double arcing which occurs when the nozzle contacts the workpiece. Such drag cups are generally made from a ceramic material or a metal material, the latter being preferred due to easy breakage of ceramic materials. When a metal material is used, an insulating layer is provided between the torch nozzle and the drag cup so as to electrically isolate the drag cup from the nozzle. In either event, the drag cup is cooled by shielding gas flowing along the outside of the nozzle during operation of the torch. Drag cups heretofore provided have included a cylindrical skirt portion providing a relatively large diameter opening, often larger than the diameter of the nozzle and providing a large open area at the outer end of the drag cup. The shielding gas flows through the open area without disturbing the plasma jet which exits through the outlet opening in the end wall of the nozzle, and the primary disadvantage of this drag cup structure is that the nozzle is not protected from the blow-back of molten metal. Accordingly, metal builds up at the outer end of the nozzle and, ultimately, will short the nozzle to the workpiece and cause double arcing. Other drag cups heretofore available have included an end wall across the outer end of the cup and provided with a small opening axially aligned with the outlet opening from the nozzle for the plasma jet. The end wall blocks the blow-back of molten metal and the latter is blown away from the end of the drag cup by the shielding gas which flows through the central opening therein for the plasma arc and through a plurality of laterally outwardly directed bleed ports provided in the end wall about the central opening. The shielding gas flow through the central opening and through the vent passages must be smooth and, accordingly, drag cups having this structure are undesirably expensive to manufacture due to the drilling and deburring of the multiplicity of holes therethrough."}
{"text": "Aromatic compounds, especially naphthalene derivatives, bearing fluorine atoms on adjacent carbons (i.e., vicinal) have been found to be useful as liquid crystal materials. They are typically made by a multi-step process, starting from the aromatic amine via a fluoro-dediazoniation process (N. Yoneda and T. Fukuhara, Tetrahedron, vol. 52, No. 1 (1996), pages 23-36).\nNo simple methods are known for producing vicinal difluoro aromatic compounds. Methods for defluorinating highly fluorinated compounds are known; but, none of the methods have been shown to produce vicinal difluoro compounds in high yield. For example:\nC. Hu, et al., Journal of Fluorine Chemistry, Vol. 48 (1990), pages 29-35, disclosed a method of synthesizing perfluoroaromatics, such as tetradecafluorobicyclo[4.4.0]dec-1 (2),6(7)-diene and perfluorotetralin, by defluorination of hexadecafluorobicyclo[4.4.0]dec-1 (6)-ene in an aprotic solvent using activated zinc powder as a catalyst. The extent of defluorination depended on the polarity of the aprotic solvent used.\nJ. Burdon and I. W. Parsons, Journal of Fluorine Chemistry, Vol. 13 (1979), pages 159-162, disclose the formation of 2,5-difluorothiophen by pyrolysis of 2,2,5,5-tetrafluoro-3-thiolen over sodium fluoride.\nSergey S. Laev and Vitalii D Shteingarts, Journal of Fluorine Chemistry, Vol. 96 (1999) pages 175-185, disclose the reductive dehalogenation of polyfluoroarenes by zinc in aqueous ammonia. In the reaction, hydrogen atoms replace fluorine atoms in the polyfluoroarenes.\nJP 2001-10995A (Ogawa, et al.) describes a four-step process for synthesis of vicinal difluoro aromatic compounds involving fluorination of a hydroxy aromatic compound to form a tetrafluoro intermediate in two steps followed by hydrogenation and defluorination under basic conditions. It also discloses reduction of a difluoroketone intermediate with aluminum isopropoxide and then base-catalyzed dehydrohalogenation to form a difluoro aromatic compound. A third method involves reaction of the difluoroketone with lithium aluminum hydride to form a fluoroepoxide, addition of HF, and elimination of water to give a vicinal difluoro aromatic compound. The best overall yield shown is  less than 50%.\nThere remains a need for an effective and simple method for preparing vicinal difluoro aromatic compounds in high yield.\nThis invention is directed to a method of preparing vicinal difluoro aromatic compounds in high yield from hydroxy aromatic compounds and to preparing intermediates thereof. The hydroxy aromatic compound can be a mono-, bi- or tricyclic aromatic in which the rings are separate or fused. One or more of the rings can contain heteroatoms, such as oxygen, nitrogen, or sulfur, and can contain substitutions, in addition to the hydroxy substitution. Substitutions on one or more of the rings can include a halogen atom, a C1 to C20 alkyl, a C5-C10 cycloalkyl, a C6 to C12 aryl, an amino, a nitro, a C1 to C10 alkyl ether or thioether, a C1 to C10 alkyl ester, a CF3, a Rxe2x80x2SO2O, \nwhere Rxe2x80x2 is CF3, a C1 to C20 alkyl, a substituted or unsubstituted C5 to C10 cycloalkyl, or a substituted or unsubstituted C6 to C12 aryl, in which the substitution on the cycloalkyl or aryl can be a C1 to C20 alkyl or a C5 to C8 cycloalkyl; Rxe2x80x3 is a C1-C10 saturated or unsaturated alkyl; x is an integer from 0 to 10, and y is an integer from 0 to 10.\nThe process can be described by the following reaction steps: \nwhere R is a hydrogen atom, a halogen atom (Cl, Br, I, F), Rxe2x80x2SO2O, CF3, a fused aryl, a C1-C20 alkyl, amino, nitro, a C1 to C10 ether or thioether, a C1 to C10 ester, a heteroaryl, wherein the heteroatom can be O, N, S, \nor R forms an aryl.\nRxe2x80x3 is CF3, a C1 to C20 alkyl, a substituted or unsubstituted C5 to C10 cycloalkyl, or a substituted or unsubstituted C6 to C12 aryl, wherein the substitution on the cycloalkyl or aryl can be a C1 to C20 alkyl or a C5 to C8 cycloalkyl; Rxe2x80x3 is a C1-C10 saturated or unsaturated alkyl; x is an integer from 0 to 10, and y is an integer from 0 to 10. The preferred R group is trans-4-propylcyclohexyl.\nThe yields obtained from reactions (1) and (2) are highly dependent on the solvents employed for these steps. Polar aprotic solvents are desirable for the electrophilic fluorination in step (1) and dimethylformamide (DMF) is particularly preferred because it unexpectedly resulted in yields of greater than 95% difluoro ketone product. Reaction step (2) can be conducted in various solvents including aliphatic and aromatic hydrocarbons, halocarbons, ethers, etc.; however, toluene unexpectedly gives much higher yields of the tetrafluoro product compared to other organic solvents. Step (3) involves reacting the tetrafluoro compound with a reducing agent, such as metallic zinc, copper, magnesium, or a mixture thereof, to form the vicinal difluoro aromatic compound in high yields (e.g., 90% or more). This reaction is preferably carried out in buffered aqueous ammonia in the presence of an organic solvent such as tetrahydrofuran (THF), methyl tert-butyl ether, acetonitrile, ethanol, or DMF. Since the tetrafluoro compound reacts under basic conditions to form a trifluoronaphthalene by-product, the pH of the aqueous ammonia is buffered to  less than 14 by addition of an ammonium salt, particularly NH4CI. Under these conditions, the selectivity to the desired vicinal difluoronaphthalene product is significantly increased.\nThis method of preparing vicinal difluoro aromatic compounds has the following advantages over known methods:\nthe difluoro ketone and tetrafluoro intermediates do not need to be purified prior to subsequent reaction,\nthe product is produced in high selectivity,\nthe overall yield is 70% or more, and\nthe product easily can be separated and purified by known methods.\nIn the method of this invention, vicinal difluoro aromatic compounds can be prepared in three steps from hydroxy aromatic compounds by electrophilic fluorination using a fluorination reagent such as Selectfluor(copyright) reagent (1-chloromethyl-4-fluoro-1,4-diazabicyclo[2.2.2]octane bis-tetrafluoroborate) to form a difluoroketone intermediate. The difluoro ketone undergoes nucleophilic fluorination by reaction with a deoxofluorinating reagent such as Deoxo-Fluor(copyright) reagent (bis(2-methoxyethyl)-aminosulfur trifluoride) to give a tetrafluoro intermediate species. The tetrafluoro intermediate is defluorinated by a metallic reducing agent in the presence of ammonium hydroxide, preferably buffered ammonium hydroxide, to provide the desired vicinal difluoro aromatic compound in high yield. An example of the reaction chemistry is described above in the Brief Summary of the Invention.\nIn the first step, a hydroxy aromatic compound (e.g., xcex2-naphthol or substituted naphthol) is reacted with an electrophilic fluorinating agent such as Selectfluor reagent, to generate a difluoroketone intermediate. This reaction can be conducted in various solvents including nitrites such as acetonitrile (CH3CN), formamides such as dimethylformamide (DMF), CH3NO2, carboxylic acids such as acetic acid, water, and an alcohol such as methanol, ethanol, and propanol.\nThe reaction can be carried out at temperatures ranging from 0xc2x0 C. to the boiling point of the solvent.\nThe fluorinating agent can be added to a solution or suspension of the hydroxy aromatic compound in one or more portions, or dropwise as a solution. Alternatively, the hydroxy aromatic compound solution or suspension can be added to a solution or suspension of fluorinating agent.\nIn the second step, the carbonyl oxygen of the difluoroketone is replaced by two fluorine atoms using a deoxofluorinating agent such as Deoxo-Fluor reagent. The reaction is carried out by reacting the difluoroketone with the deoxofluorinating agent in an organic solvent in an anhydrous atmosphere. Solvents include alkanes such as hexane, heptane, etc.; aromatic hydrocarbons such as toluene, xylenes, etc.; haloalkanes such as methylene chloride, chloroform, etc.; ethers, such as diethyl ether, THF, etc.; and any other solvent that will not react with the fluorinating reagent.\nThe reaction temperature can range from 0xc2x0 C. to 90xc2x0 C. In carrying out the reaction, the difluoroketone can be mixed with the entire charge of the fluorinating reagent or the reagent can be added dropwise to a solution of the difluoroketone. Lewis acid catalysts such as boron trifluoride etherate (BF3.Et2O) or HF can be used to accelerate the reaction. The product obtained is usually a mixture of the desired 1,1,2,2-tetrafluoro compound and the corresponding 1,1,2,4-tetrafluoro isomer. We have found that both the yield and the isomer ratio are highly dependent on the solvent used. Toluene is unexpectedly superior to other organic solvents in producing a high yield of the desired isomer. For example, when THF was used as solvent at 60xc2x0 C., a 65% yield of 1,1,2,2-tetrafluoro-6-trans-4-propylcyclohexyl)-1,2-dihydronaphthalene was obtained (1,1,2,2-tetrafluoro/1,1,2,4-tetrafluoro=51/49), while when toluene was used as solvent at the same temperature, an 85% yield was obtained (1,1,2,2-tetrafluoro/1,1,2,4-tetrafluoro=90/10) under the same reaction conditions.\nIn the third step of the process, the mixture of tetrafluoro isomers is reductively defluorinated using a reducing agent in an aqueous ammonia solution, preferably a buffered aqueous ammonia medium. Both tetrafluoro isomers react to form the same vicinal difluoro aromatic product.\nThe reducing agent in this method can be metallic zinc or other known reducing agents, such as copper, magnesium, or mixtures thereof. The metals typically are used in powder form, but other forms should also be effective. Reducing agents that do not react rapidly with water are preferred. At least one molar equivalent of the reducing agent, based on the amount of starting compound, is needed.\nThe process of reductive defluorination is preferably carried out in buffered aqueous ammonia with an organic co-solvent, such as THF. The organic co-solvent can be used to make a solution of the tetrafluoro aromatic compound. Other co-solvents that can be used include ethers such as methyl tert-butyl ether; nitriles such as acetonitrile; alcohols such as ethanol; and amides such as DMF. The aqueous ammonia is buffered to a pH  less than 14 and preferably  less than 11 with an ammonium salt; preferably ammonium chloride. Maintaining a pH below 14, minimizes the production of unwanted by-products. For example, when the reaction medium is too basic (i.e., a pH of 14), 1,1,2,2-tetrafluoro dihydronaphthalene is converted to a 1,2,4-trifluoronaphthalene compound as shown in the reaction below: \nThe rate of this reaction has been found to be pH-dependent, and consequently the selectivity to this impurity is significantly reduced as the pH is reduced below 14, especially below 11.\nTypically 3.2 ml of ammonium hydroxide and 1.6 ml of organic solvent per mmol of starting compound are appropriate for the reaction.\nThe reaction can be carried out at temperatures ranging from 0xc2x0 C. to the boiling point of the solvent; preferably 25 to 45xc2x0 C.\nThe reaction can be run in air or more preferably under an inert gas, such as nitrogen.\nThe reaction can be monitored by methods known in the art to determine completion. For example GC or GC/MS (gas chromatography/mass spectrometry) can be used to determine when the reaction is complete. Reaction times typically range from 2-48 hours.\nThe vicinal difluoro aromatic product can be isolated from the reaction mixture by methods known in the art. For example, the product can be isolated by filtering the reducing metal, extracting the aqueous layer into an immiscible organic solvent, evaporating the solvent, and purifying the product using chromatography, distillation, and/or recrystallization.\nThe invention will be further clarified by a consideration of the following examples, which are intended to be purely exemplary of the invention."}
{"text": "This is a process and method by which a paving grade asphalt material, having normal penetration of from 40 to 300, may be reduced without sacrifice of other significant properties, to less than 25 penetration by the proportioning and blending of a very hard, high softening point, solvent extracted resin-asphaltene residuum\nThe proportioning and blending is accomplished without use of further refining processes, i.e., using steam, vacuum or solvent refining. The proportioned and blended product may then be diluted with evaporative petroleum solvents to produce a penetrative binder suitable for the stabilization of pavement surface layers.\nOn Nov. 9, 1965, U.S. Pat. No. 3,216,336 was granted to J.R. Benson on \"ASPHALT PENETRATIVE GROUND STABILIZER\", hereinafter referred to as \"BENSON PATENT\". The invention related to soil stabilizing by use of an asphaltic binder mixture. The mixture made use of a selected, single, hard asphalt, resulting from the propane extraction process used in oil refining. However, the sources of suitable hard asphalt were discovered to be somewhat limited. This suitable asphaltic material could also be used to improve the quality of regular paving asphalt cements for which there was greater demand. As a result, the ASPHALT PENETRATIVE GROUND STABILIZER, covered by the \"Benson Patent\", never obtained wide general acceptance in the United States.\nFurther, due to the crude oil shortages that occurred in the early 1970's in the United States, and the necessity of producing maximum energy products and lubricating oils from the available crude oils, many oil refineries in the United States, stopped producing not only soft normal asphalt of 40 to 300 penetration, but also the 3 to 25 penetration asphalts required for use with the Benson patent. Instead, new more efficient solvent refining methods were accepted, by which even greater proportions of lubricating oils were extracted from the crude oil, from which only extremely hard or even dry powderous residues were obtained, very different from the original solvent extracted 3 to 25 penetration described by Benson. With the new (liquid separation of oil and residuum versus the previous liquid separation processes) more efficient systems, production and availability of the original 3 to 25 penetration, low softening point solvent extracted asphalts became less and less available. As a result, the use of the Benson material (marketed under the trademark \"Peneprime\") diminished to the point of non-use and relinquishment of the trademark \"Peneprime\".\nAn important part of the Benson Patent related to the character of the hard base asphalt. Previous attempts to utilize hard asphalt materials (within 3 to 25 penetration range) were unsuccessful, due to brittleness of the treated surfaces. Rapid breakup and ravelling occurred under exposed traffic. The asphalts produced in these trials were usually air-blown asphalts, asphalts where consistency (penetration) had been reduced from soft asphalts (50 penetration and above) by heating and blowing with air. While the treatment succeeded in obtaining hard asphalts, such asphalts had high softening points and very low normal ductilities. Benson discovered, however, that when soft asphalts and distillation residuals were hardened by a propane extraction process, the asphalts so produced not only had desired penetration (in the 3 to 25 range) but these asphalts also had low softening points. Benson ascribed his successful use of such materials, in a penetrative binder, to the properties of low penetration and low softening points. In only one known case (due to the peculiarity of the crude oil from which the soft asphalt was produced) could a soft asphalt be reduced by blowing with air at high temperatures, to below 25 penetration, while still retaining low softening point. All other known materials were propane extracted. Benson further patented (U.S. Pat. No. 3,399,608 - granted Sept. 3, 1968) an emulsion which also could be successfully used as a cementing binder, consisting of the hard base (solvent extracted hard asphalt) and a small quantity of suitable distillates. In still later developments (U.S. Pat. Nos. 4,268,318--granted May 19, 1981 and 4,373,961--granted Feb. 15, 1983 both to Stone) a blend of soft asphalt and a hard solvent extracted asphalt was invented, having penetrations below 50, a high ductility, and low softening points.\nPrior to discontinuing the use of the Benson penetrative binder, a number of schemes were \"improvised\" to produce satisfactory basic asphalt to meet the requirements of the Benson discoveries. Straight reduction of soft asphalts to penetrations below 25, by use of high temperatures and steam, proved to be impractical. The cost of the required equipment for such operation could not be justified. It was apparent that in order to obtain an asphaltic material equal to the propane extracted asphalts generally used under Benson's patent, a simple process utilizing locally and immediately available normal paving grades asphalts had to be developed. Also, it was essential that such asphalts be economically reduced to penetrations below 25, while retaining high ductilities and low softening points. Further, it was apparent that all necessary components of an asphaltic material meeting the requirements of the Benson Patent must be readily and economically available in all parts of the United States, using simple, low-cost equipment.\nDuring recent years, new oil refining processes have been developed. One of these processes is the ROSE process (Residuum Oil Supercritical Extraction) developed by Kerr-McGee Corporation Another, a similar process, has been engineered by Universal Oil Products (U.O.P.), INC. for Diamond Shamrock Refining and Marketing Company at its oil refinery at Three Rivers Texas. All of these oil refineries produce a resin-asphaltene residuum that is very hard. The penetration is zero at 77.degree. F., with a softening point above 175.degree. F. In some instances, due to these and other undesirable characteristics, including such materials as ash and various metals which make the residuum frequently unacceptable for blending into normal paving grade asphalt, the residuum is sometimes regarded as a \"waste\" product, even having negative value with the product makketed at substantially less than normal paving grades asphalt prices. However, it has been found that such resin-asphaltene residuum may be proportioned and blended with normal paving grades asphalt (40 to 300 penetration--hereinafter called \"Normal paving grades asphalt\"), the resultant asphalt being proportioned to give minus 25 penetration. This material exhibits good ductility and softening points, characteristics suitable for the purpose of stabilizing absorptive granular materials. In this use, the ash and metals in the resin-asphaltene residuum pose no problem since they are completely encapsulated in the absorbed cured material in the road surface.\nWhereas only propane extracted asphalt or similar products were usable under the Benson Patent, the present invention's composite product, which uses a minor part of such residuum and a major part of normal paving grades of asphalt, results in a suitable product for all of the purposes set forth in the Benson Patent. The usual proportions of residuum and diluent are less than one-third residuum and more than two-thirds normal paving grades of asphalts. Since normal paving grades of asphalts are universally available throughout the United States, it is now possible by use of these residuums, to economically produce an asphalt cement for stabilizing of highway and road building surfaces of soil and aggregates (sand and gravel). Roadway driving surfaces can be produced at one-half the cost or less of normal asphalt pavements. In addition, the new stabilizer product can be produced without expensive and complicated refining equipment. It can be produced by simple processes of blending, involving only low-cost equipment and operated by personnel having very limited technical training. The blending process consists essentially of heating the normal paving grades asphalts to temperatures of 300.degree. F. and above, proportioning, injecting or adding molten (liquid) or dry residuum until the penetration of the resultant blend is below 25 penetration. The resultant product is blended with a naptha-kerosene solvent as used in the Benson stabilization process. The napthakerosene solvent may be added to the normal paving grades asphalt before or after the injection of the residuum.\nThe end product, a blend of normal paving grades asphalts and residuum, diluted with solvent sufficient to make a fluid blend at ambient temperature, penetrates well into earthen or slightly porous highway, road or other surface. The resultant product then cures into a hard, durable surface as described in the Benson Patent. The process of blending the composite asphalt, as described above results in an increase in the polar activity of the resultant composite micelle, permitting a more complete penetration and bonding into the highway, road or other surfaces.\nWhereas normal paving grades asphalts have a substantial percentage of oily constituents, including malthenes, the resin-asphaltene residuum is essentially an asphalt residue from which nearly all oily constituents and much of the resins have been extracted. The remaining material consists of a micelle, comprising a carbonaceous core encompassed in a hard resinous material. The micellian structure of the original asphalt is thereby retained. Electrical resistivity tests indicate the retention of the resinous material since the carbonaceous core material in the center of the micelle is more or less conductive. It is this fact that accounts for the ability of the residue to disperse and dissolve in the basic asphalt. By blending these products in the disclosed proportions, the resultant mixture has a larger percentage of asphaltenes and hence greater cementing and load bearing capacities than normal paving grades asphalts. Such characteristics provide longer life roadway surfaces with greater resistance to weathering.\nThe unexpected results and advantages of blending composite asphalts from residuums and normal paving grades asphalts according to the present invention are as follows:\n1. The resin-asphaltene residuum is soluble in, and easily dispersed into, regular paving asphalts.\n2. The blend is congenial, with no significant separation.\n3. Penetration of less than 25 can be obtained.\n4. The composite of residuum and normal paving grades asphalts has an equal or greater ductility than normal paving grades asphalts.\n5. The composite blend can be diluted with a naptha-kerosene blend without precipitation of asphaltenes.\n6. The composite asphalt, when diluted with a naptha-kerosene blend, readily penetrates earthen highway, road and other surfaces, to depths equal to or greater than the diluted propane extracted asphalt as covered by the Benson patent.\n7. The treated granular material hardens and cures into a surface suitable for traffic with good durability (due to high ductility), and has weathering characteristics superior to previously used propane extracted asphalts.\n8. The resin-asphaltene residuum can be moderately diluted with kerosene or the like and shipped in tank cars to distant locations, thereafter being diluted to required viscosities.\n9. The residuums can also be shipped in containers similar to conventional shipment of roofing asphalts. It can then be melted, blended into a composite, and further blended with the required solvent to the required viscosities.\nIn the present invention, a process has been discovered which produces a material with the desired properties for use as a penetrative binder. Wide availability of residuum, from ROSE and U. O. P., locally available paving grade asphalts, and simple low-cost processing equipment now permits products equal to or better than the previously available propane extracted asphalts used by Benson."}
{"text": "Prior to the present invention, methods for determining ions in solution included flame photometry, atomic absorption photometry and ion-specific electrodes. The use of certain compounds and compositions which selectively isolate certain ions from the sample solution has become popular in ion-specific electrodes. These substances, known as ionophores, have the capability of selectively isolating ions from their counterions thereby causing a charge separation and a corresponding change in electrical conductivity in the phase containing the ionophore. Illustrative of the ion/ionophore phenomenon include ion assays utilizing membrane electrodes, liquid/liquid partitioning and fluorescence."}
{"text": "The present invention relates to a drive system for a variable vacuum capacitor.\nVariable vacuum capacitors often comprise a first set of electrodes movably arranged with respect to a second set of electrodes. The electrodes may be arranged, for example, as concentric rings. Principally, in a variable vacuum capacitor, a first set of concentric rings may be moved into a second set of concentric rings. Due to such a movement, the area is increased of parts of the first concentric rings and of parts of the second concentric rings that lie opposite to each other. This leads to an increase of the relevant active area and therefore to an increase of the capacitance of the variable vacuum capacitor. The relative movement of the first and the second electrodes is performed by means of a mechanical drive system. Such a mechanical drive system may comprise a drive screw and a nut. The drive screw may be manufactured out of stainless steel and the nut may be manufactured out of lead bronze. Apart of the electrodes and the drive system, a variable vacuum capacitor includes further elements. These further elements may comprise, for example, a housing bottom, an insulator, a housing top and a spring bellow. A first set of electrodes may be arranged in the housing bottom. Together with the spring bellow and the drive system, a second set of electrodes may be arranged in the housing top. By means of the drive system, the second set of electrodes can be moved in a given direction. The insulator may be arranged between the housing bottom and the housing top. The vacuum of the vacuum capacitor is basically surrounded by parts of the housing top and the bellows on one side, and by parts of the housing top, the insulator and the housing bottom on the other side.\nVariable vacuum capacitors are used in a wide variety of industrial applications. Such applications comprise radio broadcasting systems, manufacturing systems for flat panel displays, manufacturing systems for semiconductors, or other industrial applications. For such applications, variable vacuum capacitors are needed capable of withstanding currents of about 50 A to 200 A, or currents of up to 900 A, or higher currents, and voltages of up to 90 kV, or higher voltages. The range of the variable capacitance of a variable vacuum capacitor may be 20 pF for the lowest value and\n5000 pF for the highest value. Depending on the application, a variable vacuum capacitor may be constructed to provide other ranges of a variable capacitance. Further requirements for variable vacuum capacitors may comprise operating conditions as, for example, a certain nominal frequency, a temperature range for running the capacitor, and other requirements.\nVariable vacuum capacitors have to be constructed for a long lasting operational life. Therefore, the parts of a variable capacitor have to be manufactured in a high quality. In recent years, the requirements regarding the speed for changing the capacitance of a variable vacuum capacitor have increased. Basically, a change of the capacitance of a variable vacuum capacitor is performed by rotating the drive screw. By means of such a rotation, the capacitance of the variable vacuum capacitor is changed accordingly. If a last change of the capacitance is needed, a high rotation speed of the rotating screw is needed equally well. In recent years, the required rotation speed has increased from about 100 rpm to 600 rpm or higher. These demands have led the manufacturers to introduce new hardening techniques for the drive screw and to introduce new lubricants. However, the results achieved by these techniques have not fully satisfied the needed requirements, especially with respect to the durability of the drive system.\nIt is desirable to propose a new drive system for a variable vacuum capacitor, which drive system does not have the drawbacks of the prior art. In particular, the new drive system should have an improved life expectancy, while increasing the allowable speed for changing between different capacitance values of the variable vacuum capacitor. Preferably, the mechanical drive system comprises a drive screw and a nut, wherein the nut is arranged in a housing of the vacuum capacitor, wherein the drive screw is screwed through the nut, wherein a first electrode is arranged on one side of the drive screw, wherein, by means of a rotation of the drive screw, the first electrode is movable relative to a second electrode.\nAccording to an aspect of the present invention the nut of the mechanical drive system of a vacuum capacitor is at least partially manufactured from plastic material. In particular, such a drive system has the advantage that a high speed in the change of the capacitance of the variable vacuum capacitor can be achieved and that the durability of the drive system can be improved significantly.\nIn an embodiment variant, the nut is substantially manufactured fully out of plastic material. In particular, such a drive system has the advantage that the number of needed parts can be minimized.\nIn another embodiment variant, the part of the nut for screwing the drive screw through the nut is manufactured out of plastic material and the part of the nut for arranging the nut in the housing of the vacuum capacitor is manufactured out of metallic material. Such an embodiment has the advantage that the physical properties of the nut can be adapted to the physical properties of the material contacting the nut when installed.\nIn a further embodiment variant, the plastic material comprises a high performance plastic material. Such an embodiment has the advantage that a high performance plastic material with well known physical properties can be used for achieving desired results with respect to durability and speed of operation of the mechanical drive system.\nIn another embodiment variant, the plastic material comprises a high performance and high temperature up to 250° Celsius polymer plastic material with a friction factor of <0.1 @ 100° Celsius. Such an embodiment variant has the advantage that well known products may be used for achieving the desired results with respect to durability and speed of operation of the mechanical drive system.\nIn a further embodiment variant, the metallic material comprises a stainless steel. Such an embodiment has the advantage that a well known material for achieving durability and speed of operation can be used.\nIn another embodiment variant, the drive screw comprises a PVD (PVD: Physical Vapor Deposition) coating or a coating of a similar process, with a Vickers hardness of 6,000 HV or higher. Such an embodiment has the advantage that durability and speed of operation can be significantly improved.\nIn FIG. 1, the reference numeral 1 refers to a variable vacuum capacitor. The variable vacuum capacitor comprises a housing and insulator. A first electrode (not shown) and a second electrode (not shown) are arranged inside the housing and the insulator of the variable vacuum capacitor. For example, the first electrode may be attached to a drive screw, which has reference numeral 2 in FIG. 1, and the second electrode may be arranged opposite the first electrode in a fixed position relative to the housing of the variable vacuum capacitor. The capacitance of the variable vacuum capacitor is adjusted by means of a change of distance between the first electrode and the second electrode. In order to increase the active area, the first electrode and the second electrode may comprise ring electrodes or helical electrodes. The change in distance between the electrodes is performed by means of a mechanical drive system. In FIG. 1, the reference numeral 3 relates to a nut. The nut 3 is arranged in a fixed position relative to the housing of the variable vacuum capacitor. For example, the nut 3 may be attached to the housing by means of several screws. The drive screw 2 is screwed through the nut. Therefore, when the screw 2 is turned in one direction or the other, the first electrode is moved with respect to the housing and equally well with respect to the second electrode, leading to an increase or a decrease of the capacitance of the variable vacuum capacitor. In FIG. 1 the reference numeral 4 refers to a cap and the reference numeral 5 refers to a ball bearing. The cap 4 may be manufactured out of aluminium, for example. The ball bearing may be a standard ball bearing. The cap 4 and the ball bearing 5 may improve the quality of the mechanical drive system. The nut 3 comprises plastic material which, in contact with the drive screw 2, provides preferred physical properties, especially with respect to durability and speed of operation of the mechanical drive system.\nIn FIG. 2, a first embodiment of the nut is shown. In FIG. 2, the reference numeral 1 refers to a plastic cylinder of the nut and the reference numeral 2 refers to a metallic cylinder of the nut. The plastic cylinder 1 is arranged inside the metallic cylinder 2. The outer diameter of the plastic cylinder 1 may be a little bit larger, for example 1/10 mm, than the inner diameter of the metallic cylinder 2. The plastic cylinder 1 may be pressed into the metallic cylinder 2. The plastic cylinder 1 comprises a screw thread adapted to house the drive screw. The metallic cylinder 2 may comprise means, for example holes for housing screws, for arranging the nut in a housing of the variable vacuum capacitor.\nIn FIG. 3, a second embodiment of the nut is shown. In FIG. 3, the reference numeral 1 refers to a plastic component of the nut and reference numeral 2 refers to a metallic part of the nut. The metallic part 2 of the nut may be adapted to protect the plastic part 1 of the nut from forces applied by screws when the nut is arranged in a housing of the variable vacuum capacitor.\nWhen the mechanical drive system is installed, a lubricant may be applied to the drive screw, such that the friction between the drive screw and the nut is reduced further. Such a lubricant may comprise synthetic oil.\nThe drive screw may comprise a PVD (PVD: Physical Vapor Deposition) coating, which coating may be manufactured by means of well known methods and procedures.\nThe plastic material may be manufactured in such a manner that the variable vacuum capacitor may be operated in a temperature range from about 0° Celsius to about 125° Celsius. For example, the plastic material may comprise a high performance and high temperature polymer plastic material, which may have the following properties: maintenance free, wear-resistant, low friction, impact-resistant, vibration-resistant, well-priced, easy to process (e.g. chipping, grinding, burnishing), adapted to be glued or to be welded, anti-adhesive. The features of the plastic material may comprise: usable in a temperature range of −50° Celsius to 140° Celsius, maximal permissible temperature of 250° Celsius, dynamic coefficient of friction of <0.1 (for dry run @ a temperature of 100° Celsius), and wear factor of <0.15 mm/100 km."}
{"text": "An automatic transmission of a vehicle commonly includes a torque converter for coupling the transmission to the vehicle's engine. The torque converter provides a fluid coupling which allows the engine to spin somewhat independently of the transmission so that the transmission can slow or stop without stalling the engine. Beneficially, the torque converter also provides a magnification of torque from the engine to the transmission when accelerating from a halt. However, at higher vehicle speeds where the speed of the transmission catches up to the speed of the engine the fluid coupling is less efficient than a rigid coupling between the transmission and the engine, and results in increased fuel consumption.\nIt has been proposed to provide a torque converter with a lockup clutch which provides a selectively operable direct connection between the engine and the transmission to eliminate slippage and to improve fuel efficiency. The lockup clutch has a driving part which is driven by the engine and a driven part which is in direct connection with the transmission. The driving part and the driven part are able to be brought selectively into frictional engagement so as to provide the direct driving connection between the engine and the transmission, and are able to be brought selectively out of frictional engagement when the direct connection is not desired, for example when the vehicle is at a standstill.\nThe applicant has determined that by providing more accurate control of the lockup clutch better fuel economy may be achieved."}
{"text": "A need exists for a method for therapeutic treatment of afflicted tissue and cells of animals that is non-invasive, easy to apply, and quick to provide relief.\nThe present embodiments meet these needs.\nThe present embodiments are detailed below with reference to the listed Figures."}
{"text": "Noroviruses are non-cultivatable human Caliciviruses that have emerged as the single most important cause of epidemic outbreaks of nonbacterial gastroenteritis (Glass et al., 2000; Hardy et al., 1999). The clinical significance of Noroviruses was under-appreciated prior to the development of sensitive molecular diagnostic assays. The cloning of the prototype genogroup I Norwalk virus (NV) genome and the production of virus-like particles (VLPs) from a recombinant Baculovirus expression system led to the development of assays that revealed widespread Norovirus infections (Jiang et al. 1990; 1992).\nNoroviruses are single-stranded, positive sense RNA viruses that contain a non-segmented RNA genome. The viral genome encodes three open reading frames, of which the latter two specify the production of the major capsid protein and a minor structural protein, respectively (Glass et al. 2000). When expressed at high levels in eukaryotic expression systems, the capsid protein of NV, and certain other Noroviruses, self-assembles into VLPs that structurally mimic native Norovirus virions. When viewed by transmission electron microscopy, the VLPs are morphologically indistinguishable from infectious virions isolated from human stool samples.\nImmune responses to Noroviruses are complex, and the correlates of protection are just now being elucidated. Human volunteer studies performed with native virus demonstrated that mucosally-derived memory immune responses provided short-term protection from infection and suggested that vaccine-mediated protection is feasible (Lindesmith et al. 2003; Parrino et al. 1997; Wyatt et al., 1974).\nAlthough Norovirus cannot be cultivated in vitro, due to the availability of VLPs and their ability to be produced in large quantities, considerable progress has been made in defining the antigenic and structural topography of the Norovirus capsid. VLPs preserve the authentic confirmation of the viral capsid protein while lacking the infectious genetic material. Consequently, VLPs mimic the functional interactions of the virus with cellular receptors, thereby eliciting an appropriate host immune response while lacking the ability to reproduce or cause infection. In conjunction with the NIH, Baylor College of Medicine studied the humoral, mucosal and cellular immune responses to NV VLPs in human volunteers in an academic, investigator-sponsored Phase I clinical trial. Orally administered VLPs were safe and immunogenic in healthy adults (Ball et al. 1999; Tacket et al. 2003). At other academic centers, preclinical experiments in animal models have demonstrated enhancement of immune responses to VLPs when administered intranasally with bacterial exotoxin adjuvants (Guerrero et al. 2001; Nicollier-Jamot et al. 2004; Periwal et al. 2003). Collectively, these data suggest that a vaccine consisting of properly formulated VLPs represents a viable strategy to immunize against Norovirus infection."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a control circuit for a DC motor driving circuit. In particular, the invention is concerned with a control circuit for ensuring safety in operation of an H-bridge circuit employed for driving and controlling a DC (Direct Current) motor, and more particularly a control circuit which is capable of preventing occurrence of a through-current flow upon turning on/off of switching elements constituting the H-bridge circuit.\n2. Description of the Related Art\nFor better understanding of the present invention, description will first be made in some detail of the background technique. FIG. 6 is a block diagram showing a hitherto known DC motor driving circuit which includes an H-bridge switching circuit for driving and controlling a DC motor, FIG. 7 is a circuit diagram showing a predriver circuit serving as a control circuit for driving and controlling the H-bridge switching circuit shown in FIG. 6, and FIG. 8 is a circuit diagram showing in detail a major portion of the predriver circuit shown in FIG. 6.\nReferring to FIG. 6, a DC motor 1 is electrically connected in the form of an H-bridge circuit together with four switching elements 3a, 3b, 3c and 3d which serve for controlling the rotation speed of the DC motor 1 by controlling the duty cycle of a pulse current supplied to the DC motor 1 as well as for reversing the direction of rotation of the DC motor 1 by changing the flow direction of the current supplied to the motor 1.\nAs can be seen, the first and second switching elements 3a and 3b are connected to a power supply source V.sub.B, while the third and fourth switching elements 3c and 3d are connected to the ground potential. The DC motor 1 has one terminal connected to the junction between the first and third switching elements 3a and 3c and the other terminal connected to the junction between the second and fourth switching elements 3b and 3d. Further, the switching elements 3a, 3b, 3c and 3d have electrodes connected to output terminals of driving circuits 5a, 5b, 5c and 5d, respectively, which in turn have input terminals connected to a control unit 9 constituted by a computer or the like via a predriver circuit 7 which also serves as a control circuit for controlling the driving circuits 5a to 5d under the command of the control unit 9.\nFIG. 7 shows a circuit configuration of the driving circuit 5a and the first switching element 3a, being understood that the other driving circuits 5b to 5d and the switching elements 3b to 3d are also implemented in substantially equivalent configurations, respectively. As can be seen in the figure, the first switching element 3a is constituted by an N-channel MOSFET (Metal Oxide Semiconductor Field Effect Transistor), while the switch driving circuit 5a is constituted by an emitter-grounded transistor 51 having a base connected to a corresponding output terminal of the predriver circuit 7 and a collector connected to the power supply source V.sub.B via a resistor 52, a transistor 53 having a base connected to the collector of the transistor 51, an emitter connected to the ground potential similarly to the transistor 51 and a collector connected to the power supply source V.sub.B via a resistor 54, and a transistor 55 having a base connected to a junction between a resistor 54 and the collector of the transistor 53, an emitter connected to an input terminal of the first switching element 3a and additionally connected to the collector of the transistor 53 via a diode 56 and a collector connected to the power supply source V.sub.B. The configurations and connection mentioned above apply similarly to the other switching elements 3b, 3c and 3d and the other switch driving circuits 5b, 5c and 5d, as mentioned above.\nIn operation, when the output of the predriver circuit 7 is at a low level, the transistor 51 is in the non-conducting state (OFF) with the transistor 53 conducting (ON), as a result of which the output level of the output transistor 55 is low. Consequently, the switching element 3a (3b, 3c or 3d) is in the non-conducting or OFF-state. On the other hand, when the output of the predriver circuit 7 becomes high, the transistor 51 is turned on while the transistor 53 is turned off, which results in that the output transistor 55 is turned off. Consequently, a high level output of the transistor 55 is applied to the switching input terminal of the switching element 3a (3b, 3c or 3d) which is thus turned on.\nBy combining appropriately the ON/OFF operations of the switching elements 3a, 3b, 3c and 3d, each controlled in the manner described above, speed control, reversing of rotational direction of the DC motor 1 as well as braking and stoppage thereof can be realized with the current supply to the DC motor 1 being correspondingly controlled.\nMore specifically, referring to FIG. 6, when the first and fourth switching elements 3a and 3d are in the OFF-state, a current flows from the power supply source V.sub.B to the ground by way of the second switching element 3b, the DC motor 1 and the third switching element 3c, whereby the DC motor 1 is caused to rotate in a forward direction.\nOn the other hand, when the first and fourth switching elements 3a and 3d are turned on with the second and third switching elements 3b and 3c being off, a current flows from the power supply source V.sub.B to the ground via the first switching element 3a, the DC motor 1 and the fourth switching element 3d, causing the DC motor 1 to rotate in the backward or reverse direction.\nFurther, when the first and second switching elements 3a and 3b are turned on with the third and fourth switching element 3c and 3d being turned off, both ends of the DC motor 1 are electrically coupled through the first and second switching elements 3a and 3b, whereby the DC motor 1 is braked. Furthermore, when the first to fourth switching elements 3a to 3d are all turned off, the DC motor 1 is caused to stop. Parenthetically, in practical applications, the second and third switching elements 3b and 3c are mutually interlocked to be turned on and off in combination.\nIn the switching circuit composed of the switching elements 3a, 3b, 3c and 3d connected in the H-bridge form as described above, a short-circuit fault may take place between the power supply source V.sub.B and the ground if the first and third switching elements 3a and 3c or the second and fourth switching elements 3b and 3d should simultaneously be turned on, which may occur in dependence on erroneous timing at which operations of these switching elements are changed over. Such being the circumstances, the predriver circuit 7 is provided and dedicated to the control of the switching elements 3a, 3b, 3c and 3d with a view to preventing the occurrence of such short-circuit fault.\nDescription will now turn to the predriver circuit 7. FIG. 8 is a circuit diagram showing an exemplary configuration of a predriver circuit known heretofore. As can be seen in the figure, the predriver circuit 7 includes first to third input terminals 70a, 70b and 70c and four NOR circuits 71a, 71b, 71c and 71d provided at the output side and connected to the switch driving circuits 5a, 5b, 5c and 5d, respectively, wherein the input sides of the four NOR circuits 71a, 71b, 71c and 71d are connected in such a manner as illustrated in FIG. 8.\nMore specifically, the first input terminal 70a is connected to one input terminal of the first and fourth NOR circuits 71a and 71d, respectively, via an amplifier 72a and at the same time connected to one input terminal of the second NOR circuit 71b via an inverter 73a. On the other hand, the second input terminal 70b is connected to another one of the input terminals of the second NOR circuit 71b and one of the input terminals of the third NOR circuit 71c, respectively, and additionally connected to another input terminal of the first NOR circuit 71a via an inverter 73b.\nFurther, the first to third terminals 70a, 70b and 70c are connected to input terminals of a flip-flop circuit 74 via amplifiers 72a, 72b and 72c and inverters 73a, 73b and 73c, respectively, wherein output terminals of the flip-flop circuits 74 are connected to other input terminals of the third and fourth NOR circuits 71c and 71d, respectively.\nAdditionally, the first and second input terminals 70a and 70b are connected to input terminals of a delay circuit 75 via the amplifiers 72a and 72b and the inverters 73a and 73b, respectively, wherein output terminals of the delay circuit 75 are connected to the other input terminals of the first to fourth NOR circuits 71a, 71b, 71c and 71d via inverters 76a, 76b, 76c and 76d, respectively,\nThe delay circuit 75 is comprised of four series connections, each of which includes a transistor 10, a constant current circuitry 11 and a capacitor 12 connected in parallel to the transistor 10, wherein the four series connections are connected in parallel with one another. In response to the state change (i.e., from ON-state to OFF-state or from OFF-state to ON-state) of the input signal supplied from the control unit 9 to the first and second input terminals 70a and 70b of the predriver circuit 7, the output levels of the NOR circuits 71a, 71b, 71c and 71d constituting the output stage of the predriver circuit 7 are caused to change, whereby the switching elements 3a, 3b, 3c and 3d are correspondingly turned on and off. In that case, the delay circuit 75 serves to delay the operations of the switching elements 3a, 3b, 3c and 3d such that the first and third switching elements 3a and 3c or the second and fourth switching elements 3b and 3d are positively prevented from being simultaneously turned on for the purpose of excluding the short-circuit fault mentioned hereinbefore. As a consequence, when the operating state of the DC motor 1 is to be changed over, the power supply to the DC motor 1 is temporarily or transiently cut off.\nFIG. 9 shows signal waveforms at various circuit points in the predriver circuit 7 shown in FIG. 8 for illustrating operation thereof on the assumption that operation state of the DC motor 1 is changed over in the sequence of stop, forward rotation, braking, backward (reverse) rotation, braking, forward rotation and backward rotation in this order. More specifically, in FIG. 9, waveforms S1 to S3 illustrate changes in the output signals (control signals) of the control unit 9 which are applied to the first to third input terminals 70b, 70a and 70c, respectively, of the predriver circuit 7, while A to H represent signal waveforms at circuit points A to H shown in FIG. 8, respectively. Further, waveforms S5a to S5d represent output signals supplied from the first to fourth NOR circuits 71a, 71b, 71c and 71d to the first to fourth switch driving circuits 5a, 5b, 5c and 5d respectively.\nThe period during which the electric power supply to the DC motor 1 is cut off is referred to as the dead time. Usually, this dead time is optimally selected in dependence on the type of the DC motor 1, the switching elements 3a, 3b, 3c and 3d and other elements employed actually. With the circuit arrangement shown in FIG. 8, the dead time may be selectively determined by varying the capacity of the capacitor 12.\nIn the predriver circuit 7 for the switching elements 3a, 3b, 3c and 3d which constitute the H-bridge switching circuit known heretofore, at least two capacitors 12 (four capacitors in the case of the illustrated example) are required to be incorporated in the delay circuit 75 for preventing the short-circuit fault from occurring in the H-bridge switching circuit. In this conjunction, it is noted that when the predriver circuit 7 is to be implemented in the form of an integrated circuit, the capacitors 12 will have to be provided externally so that the individual capacitances can be adjusted by exchanging the capacitors in conformance with the switching elements used in the H-bridge circuit as well as the DC motor. In that case, the number of terminals of the integrated circuit for electrical connection of the capacitors 12 will necessarily be increased in dependence on the number of the capacitors to be externally attached, which provides a great obstacle to miniaturization of the integrated predriver circuit. Needless to say, increase in the number of circuit elements involves a problem that high operation reliability of the whole circuit can not be ensured."}
{"text": "Commercial hardening techniques for ferrous metal components typically involve the steps of heating the components in a furnace or a bath of molten salt, such as a chloride, quenching the components in a bath of another molten salt such as a nitrate, nitrite, or mixture thereof, and thereafter causing the parts to be passed through one or more rinse baths to clean off salts which cling to the parts as a result of the heating and quenching steps. A certain amount of each of the molten salts is transferred by drag-out to the rinse bath, thus requiring occasional replenishment of the molten salts, as well as replacement of the rinse water.\nIt is desirable to recover the salts from the rinse water and return them to the heat treating baths. However, a number of problems have heretofore prevented such recycling on a straightforward and economical basis. First, it is difficult to justify economically the installation of recovery apparatus unless the quantities of reclaimed salts are very high, and the recovery process can be carried out without shut-down of the heat treating system. Secondly, it has been found that the impurities which typically find their way into heat treating baths tend to build up to unacceptable levels particularly fast when the salts are recovered from the rinse and recycled into the heat treat bath. For example, chloride salts from a heating bath are dragged into the quench bath and thence into the rinse where recycling tends to return them to the quench bath. Recycled chloride impurities quickly reach insolubility levels in the quench, and crystallize out as a sludge which must be removed from time to time. In addition, chloride tends to substantially reduce the quench efficiency of a nitrate/nitrite bath and, thus, removing the chlorides is highly desirable.\nAnother impurity which will build up quickly as a result of salt recycling is carbonate, a reaction product of decomposition of the salt itself. For example, a sodium nitrate/nitrite salt tends to form Na.sub.2 CO.sub.3, the carbon being taken out of the CO.sub.2 present in the air. The sodium carbonate is quite soluble in the rinse water but only sparingly soluble in the nitrate/nitrite quench bath and, thus forms a precipitate in the bath which impairs quenching efficiency and coats mechanical components in contact with the bath.\nIt should be understood that both the chloride and carbonate impurities, as well as many other impurities such as hard water ions, are typically present in the heat treating salt baths of prior art systems. However, the levels of concentration of such impurities are kept reasonably and acceptably low for longer periods of use by reason of the fact that impurities as well as desirable salts are continuously lost by drag-out. Lost salt, of course, is replaced by fresh salt which does not contain the impurities. Where salt is recycled, however, undesirable impurities dragged into the rinse and which are soluble in the rinse water are retained in the recycled materials and, after a short time, reach unacceptable concentrations."}
{"text": "In web offset printing, a desired image is repeatedly printed on a continuous web or substrate such as paper. In a typical printing process, the continuous web is slit in the longitudinal direction (the direction of web movement) to produce a plurality of continuous ribbons. The ribbons are aligned one on top of the other, folded longitudinally, and cut laterally to produce a plurality of multi-page, approximately page-length segments, each of which is termed a signature. The term signature also encompasses a single printed sheet that has or has not been folded. Because more than one different signature can be printed at one time, it is often desirable to separate the different signatures by transporting successive signatures in different directions or paths.\nOne way to accomplish the sorting of a single stream of signatures is to use a diverter mechanism, such as a diverter wedge, to divert successive signatures to one of two paths. Once diverted, the signatures typically are transferred to a conveyor using rotating buckets (also known in the art as fans, fan wheels, paddle fans, or rotary flywheels).\nTypically, two sets of rotating buckets assemblies are utilized, one set to deliver signatures traveling along a first path to a first conveyor and the other set to deliver signatures traveling along the second path to a second conveyor. Each set of bucket assemblies includes several individual buckets arranged at a spaced distance from one another along a common axis or shaft. Each bucket has multiple aligned blades that define pockets or slots between them for receiving signatures and transferring the signatures to the respective conveyor.\nIt is desirable to increase the operating speed of a printing press in order to increase the printed product output. However, as the rotational speed of the buckets is increased, it is more difficult to ensure the reliable operation of the buckets and to ensure that signatures are not damaged. For example, signature quality problems that can occur at higher press speeds include ink offset, dog-eared edges, and defects to both the leading and trailing edges of the signatures. These and other defects can lead to paper jams in the folder, resulting in press downtime and expense.\nWhen the signatures are not snugly held between guide belts, the signatures may flutter or, when the signatures are folded signatures, the signatures may open partially. Signature fluttering and/or opening can result in damage to the signatures in the folder. Signatures may also “free fall ” (i.e. travel without any belt guidance whatsoever) as they move through the folder, resulting in the potential for the signatures to flutter and/or opening."}
{"text": "1. Technical Field\nThe invention relates to a back-illuminated type imaging device in which light is illuminated from a back-surface side of a semiconductor substrate and charges that are generated in the semiconductor substrate in response to the light are read from a front-surface side of the semiconductor substrate to perform imaging.\n2. Description of the Related Art\nThe following back-illuminated type imaging device has been proposed. That is, light is illuminated from a back-surface side of a semiconductor substrate. The back-illuminated type imaging device stores charges that are generated in the semiconductor substrate in response to the light, into charge storage areas formed on a front-surface side of the semiconductor substrate. The back-illuminated type imaging device outputs signals corresponding to the charges stored in the charge storage areas to the outside by a CCD or CMOS circuit formed on the front-surface side of the semiconductor substrate so as to perform imaging.\nIn order to absorb almost all of visible light, the semiconductor substrate (a photoelectric conversion area) the back-illuminated type imaging device, the semiconductor substrate (a photoelectric conversion area) is required to have a thickness of about 10 μm. For this reason, when the back-illuminated type imaging device is fabricated, at first structures such as the charge storage areas and the CCD are first formed on the front surface of the thick semiconductor substrate. Thereafter, a wiring layer is thereafter formed on the structures, and a supporting substrate is bonded thereto via an adhesive layer. Then, the semiconductor substrate is etched from the back-surface side until the thickness becomes 10 μm, for example. After the etching, structures such as color filters and microlenses are formed on the back surface of the semiconductor substrate so as to be aligned with the structures formed on the front-surface side of the semiconductor substrate.\nIn the back-illuminated type imaging device formed in such a manner, it is necessary to form a pad on the back-surface side or front-surface side of the semiconductor substrate so as to be connected to the wiring layer formed on the front-surface side of the semiconductor substrate. A general image sensor that is presently in wide use has such a structure that a pad opening is provided on a light incidence side. In accordance with this structure, test equipments for testing the functions of chips in the wafer are installed in a fabrication line. For this reason, if the pad opening is formed on the front-surface side of the semiconductor substrate in the back-illuminated type imaging device, it becomes necessary to drastically change or modify the test equipments, which may increase the fabrication cost.\nTherefore, similarly to the general image sensor, various methods for providing the pad opening on the light incidence side (the back-surface side) have been proposed for the back-illuminated type imaging device. One method of them is to form a through hole from the back-surface side of the semiconductor substrate, thereby exposing from the back surface the pad connected to the wiring layer formed on the front-surface side (for example, see JP 2005-285814 A and JP 2006-19653 A). However, such a method has the following problems. That is, if the through hole is formed before the color filter or the microlens are formed on the back-surface side, a material for the color filter or the microlens may be left in the through hole and cannot be removed in subsequent steps. Or, since the substrate on which the material is to be deposited has a large step formed thereon, the thickness of the deposited material is not even, which may cause a fixed pattern noise that appears as oblique lines in an image. Moreover, even if the through hole is formed after the color filter or the microlens is formed, it is necessary to remove a photoresist that defines the through hole without causing any harm to the color filter or the microlens, which is very troublesome. In particular, as described above, in the semiconductor substrate having a thickness of 10 μm, the above problems may become conspicuous, and it becomes more difficult to perform such a removal work.\nIn light of the above, as disclosed in JP 2006-339566 A, the following method may be conceived. That is, a trench is formed in the semiconductor substrate so as to extend from the front surface of the semiconductor substrate to reach the back surface thereof; the trench is filled with a conductive material, and a pad is formed on the conductive material. It can be said that in principle, it is possible to fill the trench corresponding to the pad with the conductive material. However, usually, the pad has a size of 100 μm×100 μm; therefore, in order to form a trench having substantially the same size and fill the trench with conductive material, if a CVD process is used, it is necessary to form a layer of the conductive material to a thickness equal to or larger than 50 μm and then to remove the same thickness by an anisotropic etching process. Such a process step cannot be said to be practical as a semiconductor fabrication process.\nIf the size of the trench is decreased, it may not cause any problem in view of a fabrication process. However, in such a case, there arises a problem that the resistance of the conductive material filled in the trench may increase, which may affect the device characteristics."}
{"text": "1. Field of the Invention\nThis invention relates to an actuator unlocking apparatus and a method of releasing an actuator from magnetic attraction during power-on operation, in a disk storage apparatus. An actuator for moving a transducer head on a disk storage medium is locked by a magnetic locking apparatus to prevent the actuator and the disk storage medium from being damaged by unintentional movement of them when the disk storage apparatus is not powered. This invention relates to an improved method of and apparatus for releasing an actuator from magnetic locking apparatus at power-on operation.\n2. Description of the Prior Art\nIn a disk storage apparatus, it is necessary to prevent an actuator (a voice coil motor (VCM), etc.) for moving a transducer head on a recording medium from moving due to external force and from damaging data on a recording surface and the actuator itself, by locking the actuator when the disk storage apparatus is not powered. In a small disk storage apparatus, the magnetic attraction of a permanent magnet is used to lock the actuator.\nIn this prior art configuration, owing to the magnetic attraction of -the permanent magnet, power is not consumed except when locking is released during a power-on operation. Further, in this configuration, the dimensions of the disk storage apparatus can be reduced because there is mechanically no movable part. Since the power of the actuator is used to unlock, a dedicated electric circuit for unlocking is not needed. Still further, in the configuration, component parts are merely the permanent magnet and its attaching members. Therefore, the configuration is generally adopted because of the advantage that the number of parts is small and a cost can be reduced. The present invention relates to an unlocking apparatus and method of a disk storage apparatus which unlocks an actuator from the permanent magnet.\nIn prior art, in order to release the actuator from the permanent magnet, positional information (servo information) written between a zone where the actuator is locked (parking zone) and a zone where data is recorded (data zone) on a disk recording medium is used and a seek operation is used for releasing. That is, during a releasing operation, a position of the actuator is detected by reading the positional information and the current which is supplied to the actuator is controlled. The actuator is surely released from the permanent magnet and moved to the data zone. The situation in which the releasing cannot be completed because of insufficient movement or in which the actuator collides with a crash stop due to the excess movement of the actuator is avoided by reading positional information on the disk.\nHowever, the servo information written between the parking zone and the data zone on the disk medium may disappear due to the friction between the transducer head and the disk medium if the actuator is locked and released many times, because the parking zone is a zone on which the transducer head attached to the actuator slides when the rotational rate of the disk is below a predetermined value (contact start and stop). Once the servo information is lost, releasing cannot be completed because the positional information for the releasing cannot be obtained any more. In the prior art, there is therefore a problem in terms of reliability."}
{"text": "Conventional computer equipment and, in particular networking equipment, employ air deflectors to direct and/or deflect the generated or exhausted warm air away from the equipment (e.g., in an upwards direction due to the natural flow of heated air). Some equipment manufacturers offer air deflectors that are specific to a certain model of equipment (i.e., specific to the dimensions of the equipment and/or placement of the exhaust system) and are fastened to the equipment itself. However, not all equipment manufacturers offer an air deflector associated with their products. Generally, when air deflectors are not included with equipment, generic air deflectors (e.g., curved plastic) may be used that attach to the equipment via a magnetic coupler and/or mechanism. Such a magnetic coupler can add an undesirable magnetic field near the equipment to which it is attached. Additionally, the magnetic coupler not have sufficient magnetism to keep the air deflector attached to the equipment while enduring the force of the air exhaust. Additionally, since many data center aisles are rather narrow, magnetically coupled air deflectors jutting into the aisle are easy to dislodge if bumped.\nIt is with these issues in mind, among others, that various aspects of the present disclosure were conceived."}
{"text": "Methods for electroplating articles with metal coatings generally involve passing a current between two electrodes in a plating solution where one of the electrodes is the article to be plated. A typical acid copper plating solution includes dissolved copper, usually copper sulfate, an acid electrolyte such as sulfuric acid in an amount sufficient to impart conductivity to the bath, a source of halide, and proprietary additives to improve the uniformity of the plating and the quality of the metal deposit. Such additives include levelers, accelerators and suppressors, among others.\nElectrolytic copper plating solutions are used in a variety of industrial applications, such as decorative and anticorrosion coatings, as well as in the electronics industry, particularly for the fabrication of printed circuit boards and semiconductors. For circuit board fabrication, typically, copper is electroplated over selected portions of the surface of a printed circuit board, into blind vias and trenches and on the walls of through-holes passing between the surfaces of the circuit board base material. The exposed surfaces of blind vias, trenches and through-holes, i.e. the walls and the floor, are first made conductive, such as by electroless metallization, before copper is electroplated on surfaces of these apertures. Plated through-holes provide a conductive pathway from one board surface to the other. Vias and trenches provide conductive pathways between circuit board inner layers. For semiconductor fabrication, copper is electroplated over a surface of a wafer containing a variety of features such as vias, trenches or combinations thereof. The vias and trenches are metallized to provide conductivity between various layers of the semiconductor device.\nIt is well known in certain areas of plating, such as in electroplating of printed circuit boards (“PCBs”), that the use of levelers in the electroplating bath can be crucial in achieving a uniform metal deposit on a substrate surface. Electroplating a substrate having irregular topography can pose difficulties. During electroplating a voltage drop typically occurs within apertures in a surface, which can result in an uneven metal deposit between the surface and the apertures. Electroplating irregularities are exacerbated where the voltage drop is relatively extreme, that is, where the apertures are narrow and tall. Consequently, depositing a metal layer of substantially uniform thickness is frequently a challenging step in the manufacture of electronic devices. Leveling agents are often used in copper plating baths to provide substantially uniform, or level, copper layers in electronic devices.\nThe trend of portability combined with increased functionality of electronic devices has driven the miniaturization of PCBs. Conventional multilayer PCBs with through-hole interconnects are not always a practical solution. Alternative approaches for high density interconnects have been developed, such as sequential build up technologies, which utilize blind vias. One of the objectives in processes that use blind vias is the maximizing of via filling while minimizing thickness variation in the copper deposit between the vias and the substrate surface. This is particularly challenging when the PCB contains both through-holes and blind vias.\nLeveling agents are used in copper plating baths to level the deposit across the substrate surface and to improve the throwing power of the electroplating bath. Throwing power is defined as the ratio of the through-hole center copper deposit thickness to its thickness at the surface. Newer PCBs are being manufactured that contain both through-holes and blind vias. Current bath additives, in particular current leveling agents, do not always provide level copper deposits between the substrate surface and filled through-holes and blind vias. Via fill is characterized by the difference in height between the copper in the filled via and the surface.\nAnother problem encountered in copper plating is the formation of nodules on the copper deposit. Nodules are believed to be crystals of the metal being plated and grow out of the plated surface. Nodules may range in diameter from less than 1 micron to as large as several millimeters. Nodules are undesirable for a variety of electrical, mechanical, and cosmetic reasons. For example, nodules are readily detached and carried by cooling air flows into electronic assemblies, both within and external to electronic article housings, where they may cause short-circuit failure. Therefore, the nodules have to be removed before the plated substrates are assembled into electronic articles. Conventional methods of removing the nodules involve laser inspection of each copper plated substrate followed by manual removal of the nodules by workers using microscopes. Such conventional methods leave room for worker error and are inefficient.\nAccordingly, there remains a need in the art for leveling agents for use in copper electroplating baths for the manufacture of PCBs that provide level copper deposits while bolstering the throwing power of the bath and reducing nodules."}
{"text": "1. Field of the Invention\nThe present invention relates to a dispensing apparatus and more particular to a dispensing apparatus having a dispensing assembly and a slidable locking cap for restraining the dispensing assembly in a closed position.\n2Description of the Related Art\nDispensing devices, especially those for lawn care products such as a combination of lawn fertilizer and weed control, are known, as exemplified by U.S. Pat. No. 5,996,858. Lawn products are often particulate or granular, but may be liquid, and these are spread by a consumer holding the device in an open condition while walking at a prescribed rate along a predetermined path. These devices have rotatable spouts fitted to a plastic molded container having an integral handle. Some have covers, other breakaway tabs and some have both.\nThe lawn product must be handled in a recommended manner and the dispensing devices ought to be tightly closed and kept in a safe place. Cost, ease of use and reliable structure are also of concern when designing such dispensing devices."}
{"text": "Wireless Body Sensor Nodes (WBSNs) are miniaturized, wearable systems able to measure and wirelessly transmit biological signals of patients. A major field of application of WBSNs is the ambulatory acquisition of electrocardiograms (ECGs), which evaluates the electrical activity of the heart. In this context, these devices allow long time monitoring of subjects producing little discomfort and requiring minimal medical supervision.\nA recent trend in WBSNs, driven by the progress in semiconductor technologies, has been the emergence of “smart” wireless nodes [1]. These smart WBSNs (FIG. 1) can, in addition to acquiring and transmitting data wirelessly, perform advanced digital signal processing filtering noise typically corrupting the signals and/or executing an automated diagnosis.\nIn the field of ambulatory electrocardiography, an important early diagnosis step is to separate normal and pathological heartbeats. The benefit of this separation is two-fold: first, it can provide helpful information for speeding up the visual inspections of lengthy recordings by medical staff; second, it can be used to activate a more detailed analysis of only those beats presenting pathological characteristics.\nThe second scenario can lead to two non-obvious, yet substantial, benefits. On the one hand, if a detailed diagnosis is performed off-node, it can be desirable to transmit or store only pathological beats on the WBSN, greatly reducing either the energy employed for wireless transmission or the data storage requirements, respectively. On the other hand, even if the analysis of beats is executed on the wireless node, computation effort can be reduced by activating it only when abnormal beats are detected, increasing energy efficiency.\nWhile off-line algorithms have been proposed to classify the different heartbeat morphologies [2][3], their real-time implementation poses a considerable challenge on WBSN platforms, due to their high computational requirements.\nA neuro-fuzzy classifier (NFC) approach [4] is instead potentially a good candidate for this application. Its simple feed-forward structure makes it eligible to be optimized for, and executed on, the constrained resources typically present on WBSNs.\nAn important factor to consider is the high dimensionality of the heartbeat representation problem as, for every beat, tens of samples before and after the R peak of the QRS complex peak have to be considered to perform a reliable classification.\nTo effectively address this problem, in the method according to the invention enables to classify ECG beats using a method based on random projections (RPs) that can lead to reduction in input size. The approximation error introduced by random projections is theoretically bounded, nonetheless empirical evidence shows that certain projections perform better than others. Our experiments show that even a rather simple optimization, such as the one performed by a genetic algorithm in few generations, can find a proper projection to obtain optimal classification results. In one embodiment of the invention, the method includes the RP-classifier, a filtering stage and a peak detector, used to isolate beats. In one embodiment, three-lead delineation is used instead adopted as an example of detailed analysis, and is activated by the classifier only for beats identified as pathological.\nTest heartbeats were retrieved from the MIT-BIH Arrhythmia database [5], considering all beats presenting three different morphologies: normal sinus rhythms, left bundle branch blocks and premature ventricular contractions. Experimental results show that the proposed methodology can identify more than 97% of abnormal beats, while using a small fraction of the available SoC memory and computing resources.\nNeuro-fuzzy classifiers (NFCs) [6] have been extensively studied in the literature. Their ability to explicitly express uncertainty in classification, given by the employed fuzzy values, makes them particularly well suited to the problem of heartbeat classification, as proposed in [7].\nNFCs can be effectively trained using established methods, the most common being the gradient descent algorithm described in [6] and the scale conjugate gradient introduced in [8] and [9], which is employed in this work. The approach is both computationally simpler and presenting lower memory requirements than comparable methods, like the ones based on support vector machines [2] and linear discriminants [10]. Therefore, it is more suitable for execution in embedded WBSNs.\nSeveral state-of-the-art strategies for neuro-fuzzy classification of ECGs can be distinguished based on the methodology employed to extract the features of individual heartbeats generating the classifier input. Solutions to the features extraction problem include the ones based on Independent or Principal component analysis (ICA [11]/PCA, [3]), Discrete Wavelets Transform (DWT [12]) and Discrete Cosine Transform (DCT [4]). As opposed to the method disclosed in this invention, these methods demand a computation effort not compatible with WBSN resources.\nA different approach, based on detection of morphological features, is presented in [13]; it nevertheless requires a detailed analysis of heartbeats before classification, conversely to our goal of early identification of pathological beats."}
{"text": "Fabric care, particularly cleaning, can be divided into two basic theoretical approaches, namely, water-based cleaning or dry/powder cleaning. Dry cleaning can be perceived by customers as raising certain concerns regarding solvent residue and volatile organic content (VOC), however, water-based cleaning also has disadvantages.\nFor one example, in a conventional water cleaning process, significant amounts of clean water are required for both washing and rinsing cycles. Typically, the weight ratio of consumed water to fabric is around 12:1 to 20:1. As can be appreciated, this is a challenge for clean water supply infrastructure, especially those facing water shortage issues. Moreover, many forecast water-saving as a mega-trend in sustainability. Accordingly, consumers desire to conserve water is likely to increase over time.\nThus, it is desirable to provide a fabric care composition that achieves excellent washing performance with water, whilst requiring less water than conventional water-based cleaning processes. It is also desirable to provide an efficient and environmentally friendly cleaning method compared to conventional cleaning processes."}
{"text": "1. Field of the Invention\nThe proposed invention relates to the field of pipeline transportation of liquid hydrocarbons, in particular to methods for decreasing their hydrodynamic resistance or method of drag reduction.\n2. Description of Related Art Including Information Disclosed Under 37 CFR 1.97 and 37 CFR 1.98.\nRecently, polymeric anti-turbulent additives (ATAs) or polymeric drag reducing additives (DRAs) are used for increasing capacity of oil and products pipelines. They represent a polymer solution or a polymer suspension in a liquid medium. A polymer should be soluble in a transported liquid and have a high molecular weight. Another pre-requisite for the effect of drag reduction (Toms Effect) is the turbulent mode of a hydrocarbon liquid flowing in a pipeline.\nThe introduction of a DRA into a crude oil flow in an amount from 10 to 50 grams per one ton results in increasing the pipeline capacity for 15-25%. The higher a polymer molecular weight, the lower is its concentration that is required for achieving said value of drag reduction.\nUltra-high molecular (co)polymers of higher α-olefins, which are synthesized on Ziegler-Natta catalysts, are most commonly used as a polymeric component of DRA. As to their price-quality relation, they have no match among other oil-soluble polymers. Monomers having from 6 to 16 carbon atoms are used.\nInitially, DRAs were produced as a polymer solution in gasoline (kerosene). However, injection-related problems due to its high viscosity, especially in winter, caused the present use of suspension-type DRAs. A suspension is produced, as a rule, by comminuting a product of block-polymerization of higher α-olefins, which is a rubber-like material, at a temperature below its glass-transition point, and thus-obtained polymer particles are mixed with a liquid medium that does not dissolve the polymer. The liquid medium is selected so that its density does not greatly differ from that of the polymer, in order to avoid irreversible suspension layering during storage. The polymer content of a suspension-type DRA may be up to 25% and even more, which exceeds the polymer content of a solution-type additive greatly. The DRA solid phase may comprise surfactants or other anti-agglomerator agents that prevent polymer particles from agglomerating as well as additives for inhibiting polymer oxidative destruction.\nPolymerization in a solution medium has been replaced by block polymerization also from the point of a polymer molecular weight that is significantly higher at block polymerization; consequently, a polymer thus produced is of higher quality.\nThe conventional technology of a suspension agent for decreasing hydrodynamic resistance may be divided into three steps:\n1. Block polymerization of a monomer (co-monomers)\n2. Comminution of a rubber-like co-polymer\n3. Production of a suspension stable against layering (sedimentation)\nMany patents relating to methods of producing suspension-type DRAs describe cryogenic comminution of polymers (U.S. Pat. No. 4,826,728, U.S. Pat. No. 4,720,397, U.S. Pat. No. 4,340,076). Since polymers and co-polymers of higher α-olefins are rubber-like materials with low glass-transition temperatures (for example, poly-1-octene has glass-transition temperature below −70° C.), their mechanical comminution is carried out in the liquid nitrogen medium, i.e. below their glass-transition temperatures.\nIt is well known, however, that operations employing liquefied gases are associated with higher risks and require special expensive equipment as well as certain safety measures for personnel. Furthermore, liquid nitrogen itself is an expensive material. Therefore, recently many companies started to work on creating non-cryogenic technologies of comminuting polyalphaolefins (U.S. Pat. No. 6,946,500, U.S. Pat. No. 6,894,088, U.S. Pat. No. 7,271,205, US Patent Application No. 0276566, US Patent Application No. 0287568). Preliminary impregnation with moistening agents is used for which purpose heavy alcohols, solid and liquid separating agents (derivatives of stearic acid, higher linear alcohols) as well as special equipment (homogenizers, abrasive mills, rotor-stator type grinders) are applied.\nHere, it should be noted that a suspension-type additive must comprise a finely comminuted polymer, or the process of dissolving the polymer in an oil pipeline will be too long, which affects the resulting efficiency of a DRA. On the other hand, mechanical comminution of polymers, especially fine comminution to particles having a size of 100-300 microns, leads to partial mechanical destruction of macro molecules and poorer quality of the polymer.\nMechanical destruction can be avoided by, for example, capsulated polymerization of higher α-olefins (U.S. Pat. No. 6,126,872, U.S. Pat. No. 6,160,036, U.S. Pat. No. 4,693,321, US Patent Application No. 20030013783). It consists in that catalyst-comprising monomer droplets are enclosed in a polymeric coating and suspended in a medium that does not dissolve any capsulated component. Micro-block polymerization is carried out in this way, and a dispersed polymer is produced in one step in the synthesis process. A significant disadvantage is low productivity of this process.\nAn alternative method for producing fine dispersions, which does not affect the length of a polymer molecule, consists in precipitating a polymer from a solution by adding a precipitating agent (U.S. Pat. No. 5,733,953). This method is the closest one to the present invention and is taken as a prototype.\nSaid method consists in that a polymer is precipitated, which has a high molecular weight and is synthesized by co-polymerization of higher α-olefins with a Ziegler-Natta catalyst in a solvent medium.\nThe sequence of producing such a suspension is as follows (a citation): “a low-viscosity highly concentrated polymer suspension is produced by slowly adding a liquid not dissolving a polymer (e.g., isopropyl alcohol) to a polymer solution in a solvent (e.g., kerosene). When a non-solvent is added in a sufficient quantity, a polymer precipitates from the solution in the form of small particles. A liquid is separated from the precipitate, the latter is once again washed with the non-solvent. The resulting concentrated suspension, when introduced into a flow of a hydrocarbon liquid, dissolves quickly and causes a decrease in hydrodynamic resistance.”\nThe prototype disadvantages are:                Poor quality of the polymer component        Necessity of regenerating a big volume of solvents        Pollution of the environment.        "}
{"text": "Although humans are extremely good at classifying gender from facial images, experiments have shown that most people have difficulty in classify gender from images were the faces are hairless. Furthermore, error rates in gender classification using low resolution facial images increases almost ten fold.\nIt would represent a major commercial advantage if computers could have reliable vision capabilities for classifying a person's gender, particularly from low resolution images. Although gender classification has been investigated from both a psychological and computational perspective, relatively few learning based vision methods are known for gender classification.\nGolomb et al. in “SEXNET: A neural network identifies sex from human faces,” Advances in Neural Information Processing Systems, pp. 572–577, 1991, described a fully connected two-layer neural network to identify gender from human face images consisting of 30×30 pixel images.\nCottrell et al. in “Empath: Face, emotion, and gender recognition using holons,” Advances in Neural Information Processing Systems, pp. 564–571, 1991 also applied neural networks for face emotion and gender recognition. They reduced the dimensionality of a set of 4096×4096 images to 40×40 via an auto-encoder network. Those vectors were then given as inputs to another one layer network for training and recognition.\nBrunelli et al, in “HyperBF networks for gender classification,” Proceedings of the DARPA Image Under-standing Workshop, pp. 311–314, 1992 developed HyperBF networks for gender classification in which two competing RBF networks, one for male and the other one for female, were trained using sixteen geometric features, e.g., pupil to eyebrow separation, eyebrow thickness, and nose width, as inputs.\nInstead of using a raster scan vector of gray levels to represent face images, Wiskott et al. in “Face recognition and gender determination,” Proceedings of the International Workshop on Automatic Face and Gesture Recognition, pp. 92–97, 1995 described a system that used labeled graphs of two-dimensional views to describe faces. The nodes were labeled with jets which are a special class of local templates computed on the basis of wavelet transform, and the edges were labeled with distance vectors. They used a small set of controlled model graphs of males and females to encode the general face knowledge.\nMore recently, Gutta et al. in “Gender and ethnic classification of Face Images,” Proceedings of the IEEE International Automatic Face and Gesture Recognition, pp. 194–199, 1998 proposed a hybrid method which consists of ensemble of neural networks (RBFs) and inductive decision trees.\nUsing computer vision to determine the gender of subject faces can be used in various application areas, such as the gathering of population gender-statistics from patrons at entertainment/amusement/sports parks. Gender classification can also be used by television network viewer-rating studies. Gender specific computer vision can further have application in such fields as automated security/surveillance systems, demographic studies, safety monitoring systems, and human interfaces to computers."}
{"text": "Ion implantation is an important process in semiconductor/microelectronic manufacturing. The ion implantation process is used in integrated circuit fabrication to introduce controlled amounts of dopant ions into semiconductor wafers. An ion-source is used to generate a well-defined ion beam for a variety of ion species from a dopant gas. Ionization of the dopant gas generates the ion species which can be subsequently implanted into a given workpiece.\nCarbon has emerged as a widely used dopant in the semiconductor industry for a variety of material modification applications such as inhibiting diffusion of co-dopants or enhancing stability of the doped region. In this regard, carbon dioxide (CO2) and carbon monoxide (CO) are two commonly used dopant gas sources for carbon implantation. However, CO2 and CO are prone to accumulation of deposits along surfaces of the ion chamber. Additional deposits can form along surfaces of electrodes of the ion source apparatus. The deposits may form directly from the dopant gas or from interaction of the dopant gas and/or its ionization product with the chamber components.\nSuch deposit formation is problematic. Deposits along surfaces of the electrodes of an ion implantation system create conditions susceptible to energetic high voltage discharge. Voltage discharge results in momentary drops in the beam current, commonly referred to as “beam glitching”. Deposits on the extraction aperture degrade the beam uniformity and hence the uniformity of dopant levels in the doped region. Beam uniformity and the number of beam glitches (i.e., glitch rate) during the operation of an ion source can be key metrics for the performance of an ion implantation system, such as, for example, a ribbon beam ion implantation system as commonly known in the art.\nBased on the process sensitivity, there may be an upper threshold to the glitch rate and/or beam non-uniformity beyond which the implant process cannot operate with acceptable efficiency. In the event the ion source performance degrades beyond the upper threshold, the user must stop the implant operation and perform maintenance or replace the ion source. Such downtime results in productivity loss of the ion implantation system. Hence, it is desirable to maintain proper functioning of the ion source for extended periods of time in order to perform an efficient implant process.\nAs will be discussed, among other advantages of the present invention, an improved method and system for minimizing deposits and beam glitching during an ion implantation process is desired."}
{"text": "Often components need to have coupling male and female matched pairs. The invention described herein comprises a rigidizing latch assembly (“LAR”) which, in embodiments, is a robotically compatible aerospace attachment mechanism designed to rigidly attach avionics boxes or any orbital replaceable unit to the structure of a vehicle or other structure.\nHowever, its use is not limited to aerospace or to a male/female pair which must be robotically maneuvered."}
{"text": "Portable electronic devices and tokens have become an integrated part of the regular day to day user experience. There is a wide variety of common portable and handheld devices that users have in their possession including communication, business and entertaining devices such as cell phones, music players, digital cameras, smart cards, memory token and variety of possible combinations of the aforementioned devices and tokens. All of these devices share the commonality that consumer are accustomed to carrying them with them most of the time and to most places. This is true across the various demographics and age groups regardless of the level of the sophistication of the consumer, their age group, their technical level or background.\nThese common handheld devices offer options for expandable memory. Micro Secure Digital (microSD) is the popular interface across high-end cellphones while SD and MultiMediaCard (MMC) interfaces are also available in limited models. MicroSD is the least common denominator supported by the majority of these devices and tokens (in terms of size). In addition, adaptors are available to convert a MicroSD into MiniSD, SD, MMC and USB Although most popular MP3 player (iPOD) offer's a proprietary interface, competing designs do offer standard interfaces. Digital cameras offer mostly SD and MMC while extreme Digital (xD) is another option. Micro and Mini versions of these interfaces are also available in several models. Mini-USB is increasingly available across cellphones, digital cameras and MP3 players for synchronization with laptops."}
{"text": "The present invention relates generally to a novel technique for limiting the accumulation of surface charge on a sample which is being bombarded by a primary beam and from which charged particles are being extracted. More particularly, the present invention relates to a method and apparatus for preventing the accumulation of surface charge beyond a predetermined limit by alternating during successive periods the extraction from the bombarded surface of charge of opposite signs.\nSurface analysis of a material sample by processes such as secondary ion mass spectrometry (SIMS), Auger electron spectroscopy (AES) or their equivalents involve bombarding a sample with a primary beam and removing charged particles from the sample.\nSecondary ion mass spectrometry involves bombarding a sample material with a high energy primary beam of neutral and/or charged particles (positive or negative ions, atoms or molecules). These particles impact the surface of the sample causing the emission of secondary neutral and charged particles (e.g. electrons, positive and negative ions, atoms, and molecules) some of which are subsequently mass analyzed to determine the composition of the sample. Similarly, AES involves bombarding the surface of a sample with high energy electrons which then result in the emission of Auger electrons, which are then energy analyzed to provide information on the composition of the sample.\nIf the net flux of positive and negative charges, i.e. ions and electrons, entering and leaving the sample is not perfectly balanced the sample will tend to become charged. Sample charging can destroy primary beam characteristics, interfere with the extraction of charged particles, and shift the energy of the extracted charged particles, making analysis difficult or impossible.\nIn electrically conducting samples charge accumulation may be prevented by electrically grounding the sample. The ground assures that the net charge flux of the sample will always be zero and, consequently, no charge accumulation will occur. Preventing or limiting charge accumulation in electrically non-conducting or poorly conducting materials (such as polymers, ceramics, glasses, and biospecimens) presents a more difficult problem.\nThe charge cycles commonly encountered during SIMS analysis of an electrically nonconducting sample exemplify the conditions toward which the present invention is directed:\n(a) The most common mode of SIMS analysis is to use a primary beam of positive ions which results in emission of secondary positive and negative ions and secondary electrons from the sample. For known materials, the positive primary beam creates fewer secondary positive ions than primary ions. Use of an electric field of negative potential to extract positive secondary ions consequently causes the sample to accumulate a positive charge. Use of an electric field of positive potential to extract negative secondary ions and electrons causes the sample to accumulate a positive charge at an even greater rate.\n(b) A neutral primary beam for SIMS analysis causes charge accumulation at a rate that is significantly less than that of the positive primary beam, but the imbalance between positive ion and negative ion and electron production and extraction can still cause charge accumulation on the sample.\n(c) A negative primary ion beam has been used for SIMS analysis to create positive and negative secondary ions. Extraction of positive secondary ions will cause the sample to accumulate negative charge. Charge accumulation on the sample due to extraction of negative secondary ions will depend upon the material and bombardment energy. The secondary negative charge yield under keV energy negative ion primary bombardment is greater than the charge flux of the primary beam for many materials. Extraction of negative particles will cause such materials to accumulate positive charge.\nIn the prior art, conducting grids and electron flooding have been used to reduce accumulation of positive charge in non-conducting samples during SIMS analysis.\nU.S. Pat. No. 4,680,467 issued July 14, 1987, to Bryson III et. al. discloses the use of an electrically conductive grid comprised of intersecting wires or vanes positioned in closely spaced relationship to the surface of the sample for smoothing the gradients in the electrical potential in the region of the beam spot during electron spectroscopy. However, conducting grids are known to add a background signal to the mass spectrum of the sample. In addition, conducting grids have a high probability of introducing contaminants to the surface of the sample and may cause chemical changes in the surface of the sample.\nU.S. Pat. No. 4,818,872 issued Apr. 4, 1989, to Parker et al. discloses an electron gun which floods the sample with low energy electrons neutralizing charge accumulation. Electron flood guns are effective in many cases but may damage the sample, particularly polymers. Electron stimulated desorption and resulting noise may also be caused by electron flood guns.\nTherefore, in view of the above, it is an object of the present invention to provide an apparatus and method for maintaining a neutral surface potential for a sample of electrically nonconducting or poorly conducting material that is being analyzed by a method causing extraction of charged particles from the sample.\nIt is another object of the present invention to provide a non-intrusive apparatus and method for maintaining a neutral surface potential so that sample integrity is maintained.\nA yet further object of this invention is to present a method for limiting surface charge accumulation within predetermined values.\nIt is another important object of this invention to present an apparatus and method which enables the near simultaneous collection of both positive and negative secondary ion spectra.\nAdditional objects, advantages and novel features of the invention will be set forth in part in the description which follows and in part will become apparent to those skilled in the art upon examination of the following or may be learned by practice of the invention. The objects and advantages of the invention may be realized and attained by means of the instrumentalities and combinations particularly pointed out in the appended claims."}
{"text": "The present invention relates to thermally stable polycrystalline diamond compacts and more particularly to brazing such compacts to cemented carbide supports and to themselves.\nWell known in the super abrasive art are compacts of polycrystalline abrasive particles typified by polycrystalline diamond and polycrystalline cubic boron nitride (CBN) compacts. Such compacts are represented by U.S. Pat. Nos. 3,745,623 and 3,608,818 with respect to polycrystalline diamond compacts and U.S. Pat. Nos. 3,767,371 and 3,743,489 with respect to polycrystalline CBN compacts. While such polycrystalline compacts represent a significant contribution to the art in many fields of use, thermal degradation at an elevated temperature, e.g. above about 700.degree. C., did limit their usefulness, especially in metal matrix bond applications. Thermal stability of such polycrystallline compacts was improved with the advent of thermally-stable porous self-bonded diamond and CBN compacts containing less than about 3 percent non-diamond phase, hereinafter termed \"porous compacts\". Compacts of this type are the subject of U.S. Pat. Nos. 4,224,380 and 4,288,248.\nEuropean Patent Publication No. 116,403 describes a thermally-stable diamond compact comprising a mass of diamond particles present in an amount of 80 to 90% by volume of the body and a second phase present in an amount of 10 to 20% by volume of the body, the mass of diamond particles contain substantially diamond-to-diamond bonding to form a coherent skeletal mass and the second phase containing nickel and silicon, the nickel being in the form of nickel and/or nickel silicide and the silicon being in the form of silicon, silicon carbide, and/or nickel silicide. British patent application No. 8508295 describes a thermally stable diamond compact comprising a mass of diamond particles present in an amount of 80 to 90% by volume of the compact and a second phase present in an amount of 10 to 20% by volume of the insert, the mass of diamond particles containing substantially diamond-to-diamond bonding to form a coherent skeletal mass and a second phase consisting essentially of silicon, the silicon being in the form of silicon and/or silicon carbide.\nDiamond, synthetic or natural, is very difficult to wet, making the attachment of diamond to a variety of substrates difficult. Since porous compacts essentially are composed only of diamond, they are difficult to bond to cemented carbide supports, for example. In order to successfully bond porous compacts to cemented carbide supports, a bonding agent that is able to wet and/or adhere to both the diamond surface and the support is required. The bond between the diamond and carbide support must be sufficiently strong so as to provide a useful implement. Temperature conditions required to form the bond must be below those that cause graphitization of the diamond and/or damage to the carbide support system. The bonding agent also must be chemically and physically compatible during and after the bonding operation with both the porous diamond compact and carbide support.\nThe formation of a bond between porous compacts and carbide supports is inhibited by gaseous/liquid/solid films. The inhibiting films can be non-carbon material as well as adsorbed species including water vapor or oxygen. The inhibiting films, by preventing the wetting of the compact and/or carbide support by the bonding agent, reduces the strength of adhesion. Further, blocking of pores in the porous compact by impurities prevents penetration of the bonding agent into the core and reduces the ultimate bond strength.\nA wide variety of bonding agents, or brazing alloys, have been proposed in the art. With respect to thermally stable diamond compacts, British Pat. No. 2,163,144 reognizes the difficulties encountered in brazing such compacts and, thus, proposes to bond a variety of thermally stable diamond compacts to carbide supports utilizing a silver or gold-based alloy containing between about 1 and 10 percent by weight of an active metal selected from the group of tungsten, titanium, zirconium, hafnium, vanadium, niobium, tantalum, chromium, and molybdenum. The alloy has a liquidus above 700.degree. C. Data, however, only is given for diamond compacts containing a second phase of silicon, nickel, or their carbides. When the present inventors utilized brazing alloys containing carbide-forming active metals as proposed in this British patent, but when bonding porous compacts (as defined herein), bonding was erratic. Moreover, relatively low bond strengths were realized even when bonding did occur. U.S. Pat. Nos. 3,894,673 and 4,018,576 propose to bond signle crystals of diamond to metal using nickel/chromium or cobalt/chromium braze alloys."}
{"text": "This invention relates to chip pads, and in particular to circuitry associated with chip pads.\nFIG. 1 is a schematic diagram of an integrated circuit (IC). The IC is formed as a set of electrical components on a semiconductor die 1. The components include a main processing section of active components 2. Around the periphery of the die is a set of bonding pads 3, known collectively as a pad ring. The bonding pads serve to allow the die to be connected to other devices. For example, when the die is packaged, wire bond connections may be made from the pads 3 to solder pads on the exterior of the package. Those solder pads may then be soldered or clamped to a circuit board.\nEach bonding pad 3 has interface circuitry 4 located between it and the active components 2. FIG. 2 shows a simple example of the form the interface circuitry may take. In this example the interface circuitry 4 provides electrostatic discharge (ESD) protection and level shifting. Diodes 5, 6 are arranged between the pad 3 and voltage supply rail 7 and ground 8. They resist ESD interference being passed to the line 10 that connects the pad to the active components 2. A level shifter 9 converts between the voltage level used to represent “high” in the domain of the active circuitry 2 and the voltage level used to represent “high” in the domain of the pads. In this example the voltage level used to represent “high” in the domain of the active circuitry is slightly lower than the voltage level used to represent “high” in the domain of the pads, so for output signals the level shifter converts a “high” signal input from the active components 2 into a slightly higher voltage that will be passed to the pad 3, and for input signals the level shifter converts a “high” signal input from the pad 3 into a slightly lower voltage that will be passed to the active components 2.\nOne trend in modern IC design is to drastically reduce the power consumption of ICs. One way to reduce power is to divide the IC into a set of power domains. A power domain is a set of components on the IC that can together be placed in a low power mode independently of the other components on the IC. Dividing the IC into multiple power domains allows parts of the IC that are not being used to be powered down, so that they are using little or no power, whilst other parts of the IC continue to function. Normally the IC will be equipped with a power controller which determines which power domains need to be turned on to satisfy the current requirements on the IC, and arranges for the other power domains on the IC to be powered down.\nThe IC might be intended to generate outputs that will be read by other devices. Those outputs will be presented as high or low voltages, representing binary ones or zeros, at its output pads. The state of an output pad will be determined by active circuitry on the IC. That active circuitry will generate a logical output, and the appropriate output pad will be set to high or low accordingly.\nOnce the active circuitry has generated the logical output, e.g. at point 10 in FIGS. 1 and 2, the active circuitry has completed its task. Assuming no more work is required of the active circuitry, the active circuitry could in principle be powered down to reduce energy consumption. However, if the chip pad has a simple driver circuit as shown in FIG. 2, if the active circuitry is powered down then the pad will no longer carry an output, or could revert to a default state. It is generally not known when the external device that is intended to consume the output has actually read the state of the pad. Therefore, it is normally desirable to keep the active circuitry powered up for some time after it has completed its processing, just to maintain an output at a pad that can at some stage be accessed by a consumer of that output. Keeping the active circuitry alive in that way consumes power that could otherwise be saved.\nOne way to address this is to make the entire pad ring, including those latch circuits, a separate power domain from the active circuitry, and to provide a buffer as part of that power domain into which the pad state can be written when the active circuitry is to be shut down. This provides advantages over keeping the active circuitry powered up, but the requirement to store the pad state into the buffer increases the time taken to power down the active circuitry, and it can be complex to lay out the IC to that the buffer can communicate with the pads and share their power domain."}
{"text": "Combustion control of engines has long been used to achieve various purposes. The three most notable purposes are to maximize power output, to maximize fuel economy, and to minimize emissions. In particular, emissions control has become a major factor in the development of engine controls.\nVarious conventional techniques exist for controlling the emissions of an engine. Most techniques deal in some form of controlling the amount of fuel, the amount of air or other combustible gas, and the timing of the ignition. These techniques vary significantly in the precision with which they are able to control emissions to a desired level. U.S. Pat. No. 4,867,127 to Quirchmayr, et. al. discloses regulating the combustion mixture pressure to achieve a desired NOx level. While this method may present some advantages over past strategies, further improvements are still sought by the industry."}
{"text": "The present invention relates to derivatives of 3,5,3'-triiodothyronine (T.sub.3) which are useful for the measurement of free T.sub.3 in a liquid sample in which the T.sub.3 is present in both free and bound states.\nT.sub.3 is one of the two physiologically important iodinated hormones, the other one being thyroxine (T.sub.4). In serum, most T.sub.3 is bound to thyroxine-binding globulin (TBG), although some binding to albumin and probably thyroxine-binding prealbumin (TBPA) also occurs. Of the T.sub.3 present in serum, approximately 0.2 to 0.4 percent remains free; i.e., unbound to plasma proteins. It is this free fraction that is presumably active in stimulating metabolism. Moreover, it has been estimated that 65 to 75 percent of the total metabolic effects in man is due to T.sub.3. Consequently, the measurement of free T.sub.3 in blood serum or plasma is of clinical interest for the prophylaxis or treatment of a variety of disorders and diseases."}
{"text": "Side dumping vehicles for a variety of materials are known in the prior art and some examples of the patented prior art are contained in the following U.S. Pat. Nos. 3,083,058, 3,712,675, 3,205,011, 3,844,617, 3,937,502.\nThe present invention is concerned with a broadly similar class of towed side dumping vehicle which is particularly constructed for transporting picked or stripped cotton from a harvester to a module. The objective behind providing the particular specialized cotton transporting and dumping vehicle is to enable the cotton farmer to keep the harvester working in the field and avoid the high cost and wasted time of using the harvester to deliver picked cotton to the module and then deadheading back to the end of a row. By means of the invention, the expensive harvester can be kept working more hours per day anc can take advantage of good harvesting weather, thus reducing the risk of crop deterioration in bad weather.\nWith the above in mind, the invention has been tailored to meet the problems which are peculiar to handling, transporting and dumping the fluffy lightweight cotton bolls, as distinguished from other materials having vastly different physical properties. The dump body or basket of the vehicle is constructed with wire mesh panels strategically located in its bottom and opposite end walls to permit the escape of air when the picked cotton is delivered into the dumping basket from the harvester. Without these air escape panels in proper locations, the loose cotton would simply boil and float over the open top of the basket on the cushion of air built up therein under the cotton. No known prior art vehicle is suitable in this respect for handling picked or stripped cotton.\nA further feature of the invention adapting it to its very specialized mission resides in the provision of a dump body or basket having a capacity to hold two to three bales of picked cotton or two and one-half bales of stripped cotton. To achieve this, the basket must have a load holding volume of about 850 cubic feet.\nUnder such capacity or loading, the stability of the transport and dumping vehicle becomes of prime importance, and herein lies a further important aspect of the invention distinguishing it from the generally similar devices in the prior art. The geometric arrangement of the elevated dump basket pivot relative to the main frame wheels and the location of the dump basket hydraulic cylinders in relation to the main horizontal wheeled frame are quite critical to achieving stability particularly during dumping and retaining within reasonable limits the force required to raise and dump the loaded basket. A unique geometric arrangement of critical pivots and planes have been created in the invention to render the invention safe and very efficient as well as simplified and economical to manufacture.\nOther features and advantages of the invention will become apparent during the course of the following description."}
{"text": "Some types of sensors, such as angular position sensors for example, are configured to generate two or more sensing signals that vary sinusoidally with respect to a sensed parameter value. The two or more sensing signals may include a first sensing signal that varies with respect to the parameter value according to a sine function and a second sensing signal that varies with respect to the parameter value according to a cosine function. A decoder may be used to map the values of the two sensing signals to an output signal that corresponds to arctangent of the first sensing signal divided by the second sensing signal. Conventional decoders may perform the arctangent function digitally by using a look-up table."}
{"text": "In means of transportation, in particular airplanes, storable tables are frequently provided for the passengers. In VIP airplanes and in the first class section of conventional commercial passenger airplanes, a table of this type is usually pulled out of what is known as a credenza, which is arranged next to the seats, and is folded into a use position, in which the table top is arranged partially over the seat surface of the seats.\nIt is problematical in tables of this type that a plurality of requirements have to be united which are difficult to combine with one another. The credenza is usually lower than the armrests of the seats. This means that the table top has to be arranged considerably above the upper edge of the credenza in the use position. A mechanism therefore has to be provided which holds said table top in this higher position. The table top has to be self-supporting, and all loads which act on the table top therefore have to be conducted into the credenza. This can be problematical, in particular, in a long lever arm for the active forces (for example, a passenger who stumbles along the aisle).\nIt is known to secure the table top of a table of this type pivotably on a carrying arm which protrudes out of the upper edge of the credenza in the use position and therefore produces the necessary spacing between the upper edge of the credenza and the table top. The lid of the credenza is configured as a type of catching hook, and that edge of the table top which points away from the passenger seat can be hooked into the lid in its open position, with the result that the tilting moment which acts on the table top is conducted into the credenza via this lid. This solution is disadvantageous, as a lid which is fastened to the credenza by means of a hinge is not particularly well suited for absorbing loads."}
{"text": "Generally, a fixing device mounted in an image forming apparatus employs a lamp that emits infrared rays, e.g., a halogen lamp, as a heat source, or employs a method of heating with Joule heat by electromagnetic induction.\nIn general, fixing devices include a heating roller or a fixing belt coupled to a plurality of rollers and a press roller. It is necessary to reduce the heat capacity of each component as much as possible and to further concentrate heating areas of the fixing device in order to maximize overall thermal efficiency.\nFurthermore, typical heating widths are wide, thus making it difficult to intensively provide widely distributed thermal energy to a nip portion.\nIn addition, fixing quality within the fixing device for an electronic photograph is affected if heat generated unevenness is present in a paper transport direction and in a perpendicular direction thereto. Particularly, with respect to color printing, heating unevenness produces differences in color developing and can generate a gloss-like appearance.\nAdditionally, in a fixing device in which heat capacity is extremely reduced, additional problems with respect to speed irregularity, warp or deterioration of the belt, and/or expansion of the transport roller exist in that the temperature of parts of the fixing device through which paper does not pass is increased dramatically. Due to energy conservation concerns, heating of such areas is not preferable. As such, due to environmental concerns, an apparatus or method that provides energy to the nip portion which heats only the passing area of the paper or the image forming area in the paper has become an area of focus in the field.\nFurthermore, the area adjacent to the image forming area (referred to as an intermediate area) and the area inside the lateral width of the paper is in contact with the image forming area. As such, the heating of such areas affects the heating of the image forming area and the like. It is preferable that the intermediate heating area should have a low temperature in order to conserve energy, but the paper is heated during printing, and the paper typically expands and contracts due to heat. Differences in temperature at the border of the paper result in wrinkling. In addition, a diameter of the press roller (for pressing the paper from a rear surface thereof) changes depending on the temperature. Accordingly, the border between the intermediate area and the image forming area has an uneven shape thus reducing the overall print quality of the paper."}
{"text": "The present invention relates generally to a substrate polishing apparatus for chemical mechanical polishing of a semiconductor substrate. More specifically, the present invention is directed to a substrate supporting pad, to the various constructions and compositions of such a pad, and to polishing apparatus using such a pad to support a substrate during polishing.\nChemical mechanical polishing (CMP) is a process performed by a substrate polishing apparatus for planarizing or polishing semiconductor substrates, also known as substrates, to provide substantially planar front and/or backsides thereon. The CMP process is used to remove high elevation features, or other discontinuities, which are created during the fabrication of semiconductors on the substrate.\nThe CMP process is performed by a substrate polishing apparatus having a carrier head assembly that removably retains a substrate to be polished, the carrier head urging the substrate against a rotating polishing pad. Further, the carrier head assembly may or may not rotate and translate the substrate relative to the rotating polishing pad. As the polishing pad rotates, it tends to pull the substrate from beneath the carrier head assembly. To eliminate this problem, the carrier head assembly includes a substrate confining ring that extends circumferentially around the lateral edge of the substrate, and retains the substrate and limits movement of the substrate relative to the carrier head assembly. In addition, a controllable down force is applied on the substrate to urge it against the polishing pad.\nA down force is applied according to U.S. Pat. No. 5,795,215 that discloses a carrier head assembly with a vertically expandable bellows chamber wherein fluid, preferably air, is pumped into and out of a primary pressure chamber to control the load to be applied to a substrate. When air is pumped into the primary pressure chamber, the pressure in the chamber increases and a base assembly is pushed downwardly. A substrate backing member is attached by a vertically expandable bellows. The bellows and the substrate backing member may be formed of stainless steel.\nA down force is applied according to U.S. Pat. No. 5,449,316 that discloses a carrier head assembly for polishing a substrate by providing a downwardly opening plenum covered by a flexible membrane. When pressurized fluid is introduced to the plenum, the membrane applies a uniform downward pressure across a backside of a substrate to be polished. The membrane is several hundred microns in thickness and is composed of a synthetic rubber.\nA down force is applied to a substrate according to U.S. Pat. No. 5,931,719 that discloses an inflatable bladder located beneath the polishing pad. The inflated bladder is used to vary the pressure exerted against the bottom of the polishing pad, which results in varying the down force of a substrate that is held against the top surface of the polishing pad. Varying the pressure is said to compensate for center slow polishing arising when a polishing pad has been used repeatedly for polishing multiple substrates.\nUndesired overpolishing of the substrate can be attributed to a warped or otherwise unevenly planar substrate, and by the manner in which the substrate is retained to the carrier head assembly. The carrier head assembly includes a generally planar lower carrier face or platen on which is mounted a conformable and resilient carrier pad film in the form of a thick, resilient and solid film that contacts the backside of the substrate. The conformable and resilient carrier pad film may be a poromeric carrier pad film, commercially available from Rodel, Inc., Newark, Del. USA, and known as xe2x80x9cDF-200xe2x80x9d and xe2x80x9cR200T3xe2x80x9d or xe2x80x9cR200T4-470-510xe2x80x9d. Such a carrier pad film is a dense, porous, closed-cell polymer foam having a thickness generally less than 0.030 in. thick. When wetted, a fluid surface tension is provided for adhesion of a substrate in contact with the carrier pad film. Alternatively, a conformable carrier pad film may be provided by a wax mound against which the substrate is pressed to form a conformable film surface to receive a substrate for polishing. The wax material has limited resiliency.\nThe lower carrier face may lack a desired planar orientation, and may include protrusions or be misaligned from planar, all of which are irregularities that contribute to uneven distribution of a down force against the carrier pad film, and, in turn, against the backside of the substrate during polishing. Such uneven distribution will exert localized higher down force distribution on a portion of the substrate, which contributes to uneven rates of polishing and deviation from a planar polished surface on the substrate. Further, the substrate may become over-polished at a locations of localized higher down force distribution, which reduces the yield of useable planar semiconductor areas on the substrate. A conformable and resilient carrier pad film provides only limited improvement in attaining an even distribution of the applied down force. For example, the load distributing properties of the carrier pad are limited, and are often inadequate to compensate or overcome misalignment of the lower carrier face itself, and misalignments of the carrier pad film and the substrate on the lower carrier face. The carrier pad film itself may have variations in its thickness and constituent construction that might contribute to nonuniform distribution of down force. Further, the compressibility of a carrier pad film is limited by the amount of compression provided by its limited resiliency. A need exists for a conformable carrier pad that is more compressible than the resiliency provided by a carrier pad film, and that can be sufficiently compressible and/or resilient to compensate for uneven distribution of applied down force, as provided by a misaligned carrier face and by protrusions on the carrier face.\nA recurring problem in chemical mechanical polishing is the so-called xe2x80x9cedge effectxe2x80x9d, i.e., the tendency of the substrate edge to be polished at a different rate than the center of the substrate. The edge effect typically results in overpolishing or under-polishing in spite of the use of carrier films. Further, variations in the composition and construction of the carrier film itself can contribute to nonuniformly distributed, localized pressure zones at certain areas across the contacting surface on the carrier film and, hence, across the substrate to be supported by the contacting surface.\nThe compressibility of a resilient film carrier pad is limited, and the resultant pressure on the substrate is directly proportional to the degree of resilient compression by the carrier pad. Since in practice, the film carrier film must be changed regularly, it would be desirable to provide carrier pads of uniform construction and of uniform composition from one carrier pad to another, which enables the same level of internal pressurization.\nThe present invention provides a carrier head assembly for a substrate polishing apparatus applying uniform down force distribution against a substrate by using an internally pressurizable carrier pad. Advantageously, the internally pressurized carrier pad conforms to, and thereby compensates for, variations in substrate thickness, carrier assembly irregularities, and avoids the limits to compressibility and the nonuniform distribution of down force due to variations in composition and thickness, as provided by conformable and resilient carrier pad films.\nEmbodiments and advantages of the invention will become apparent by way of example from the following detailed description taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThis invention relates to nonionic fluorochemical surfactants useful as antistatic agents and lubricants for polymeric shapes.\n2. Description of the Prior Art\nU.S. Pat. 3,702,870 -- Pittman et al, issued Nov. 14, 1972, and U.S. Pat. 3,758,538 -- Litt et al, issued Sept. 11, 1973, describe preparation of fluorinated alcoholates by reaction of a fluoroketone with an alkali metal fluoride. In Pittman et al, the alcoholates are then reacted with acyl halides to obtain esters useful in improving water-repellency and oil-repellency of textiles. In Litt et al, the alcoholates are then reacted with olefins to obtain fluorine containing ethers useful as surfactants.\nU.S. Pat. No. 2,723,999 -- Cowen et al. issued Nov. 15, 1955, describes preparation of nonionic surfactants by reaction of fluorinated alcohol with ethylene oxide.\nAlthough these patents teach preparation of fluorinated surfactants, there is a definite need for improved fluorinated surfactants having useful properties as antistatic agents and lubricants for polymeric shapes."}
{"text": "Read Only Memories (ROMs) have become widely utilized in integrated devices today. They are used, amongst other things, to store data and programs. As the market moves towards system on chip solutions, the need for large on-chip ROMs has increased. The area that is used by these ROMs is usually between 5% and 30% of the overall chip area and can be as much as 50% of the total device area. Thus, the smaller the ROM can be made, the cost thereof will decrease.\nThe data and programs are stored in the ROM in the form of groups of 1's and 0's (binary code), or bits, known as words. The words are permanently stored, can only be read from the memory, and are typically made up of multiples of 8 bits or bytes. Bits forming words are, in turn, derived from compartments of equal area formed into rows and columns of a ROM array. Each compartment has an address.\nReference is now made to FIG. 1 which illustrates, by way of example, a 128 Kb ROM 10 in accordance with the prior art. ROM 10 comprises an array 12 containing 8 K words of 16 bits each. The bits of the words are configured as 256 rows and 512 columns within array 12. Each bit of each word is derived from a distinct physical unit or cell 13, equivalent to the compartments, imprinted on a silicon wafer utilizing digital Complementary Metal Oxide Semiconductor (CMOS) technology. A cell has a minimum manufacturable size for a particular process based on the components, e.g., transistors, that are required to be imprinted therein. Cells occur at the cross-section of rows or word lines 15 with column lines or bit lines 17 and are all the same dimensions (and area) for a given array 12. FIG. 1 shows a cell 13A containing a bit of 1 at the intersection of word line 15A with bit line 17A. The location of particular cells 13 correspond to an address on the ROM array 12.\nROM 10 further comprises an X-decoder 14, a Y-decoder 16, a selector 18 and a sense amplifier and output driver 20. X-decoder 14 decodes 8 bits of an address which is the part of the address that determines the correct word line and activates one word line out of 256. The Y-decoder 16 selects a number of columns or bit-lines, corresponding to a word. In the present example, it connects one of every 32 cells 13 in the 16 groups of 32 cells 13 in the activated word line, in accordance with a supplied address. Y-decoder 16 achieves this by driving selector 18, which is, in effect, a switch with multiple inputs. Thus, for each address the 16 bits of a word are selected.\nSense amplifier and output driver 20 then sense whether there is a current in bit-line 17 when connected, i.e., whether a transistor (not shown) exists in each addressed cell 13 (whether the bit is a 1), and outputs the sequence as 16 bit words. If there is no transistor, a zero bit is registered. The sensing amplifier 20 is connected to a logic control unit 22 which controls its operation.\nReference is now made to FIG. 2 which is an illustration of a prior art layout of the two basic options for cell 13, that is, both with ad without a transistor. FIG. 3 shows an exemplary layout of four by four ROM cells 13 in an array 12. In FIG. 3, transistors or non-transistors connected to word lines 15, bit lines 17 and a ground connection 41, are shown.\nThe basic cell 13 contains one transistor 24 or a non-transistor 25 which is made up of the constituent parts of a transistor but not connected to form a transistor. Transistor 24 is formed so that the word line 15 (FIGS. 1, 3) is connected to its gate 30. The drain 32 and source 33 of transistor 24 (with transistor) are connected to the bit line 17 (FIG. 3) and a shared ground 41 (FIG. 3) respectively. The equivalent areas of non-transistor 25 are likewise connected but do not conduct current.\nTransistor 24 is composed of a diffusion layer 34 of width W, a polysilicon or poly layer 36 of length, L and a contact 38 overlapped by a metal-1 layer 40. The diffusion layer 34, where it touches with poly 30 at gate 36, forms the active area or channel of transistor 24 and is a measure of the size of transistor 24 in terms of the current it draws. A non-transistor 25 does not have the diffusion layer 34 close to poly layer 36. There is therefore no active area or channel under poly layer 36 in this cell 13.\nThe minimum width and length, W.sub.min and L.sub.min of the transistor cell 24 is determined by the minimum allowed transistor size for a specific process. The minimum area of a cell 13, in general, is a function of a number of factors. One of these factors is the overlap of the diffusion layer 34 over the contact width. Therefore, an increase in the width W of diffusion layer 34 may affect the area of cell 13 if it increases this overlap beyond a certain limit. Another factor influencing the minimum area of a cell, 13 is the length, L, of the poly layer 36.\nFIG. 3, which is referred to hereinabove and illustrates an exemplary prior art layout of four by four ROM cells 13 in an array 12, is now further referred to. The rows of memory cells 13, comprising non-transistors 25 and transistors 24 share the same word line 15 as shown. Each transistor 24 or non-transistor 25, of a row is shown connected to a different column or bit line 17. When a word line 15 is activated, only those cells in that row which contain a transistor 24 will conduct current from their bit line 17 to group 41 (or to a virtual ground or to a reference). Which of the cells has a transistor 24 can be determined by sensing the currents on the bit lines 17. The needed data is conveyed by pre-programming the location of the transistors in the array. As is illustrated, two basic cells 13 share the same contact 38 and all the contacts are connected together by the metal line of bit line 17. The shared ground line 41 formed by the meeting of diffusion layers 34 from two transistors 24 or non-transistors 25 is connected to ground by a metal line (not shown), for example, every 16 cells in order to save space."}
{"text": "Integrated semiconductor circuits (ICs) can be damaged by transient pulses or overvoltages, which are coupled in via terminals (pads) or directly in lines, such that they become nonfunctional or are even destroyed. Such pulses or overvoltages may occur for example in the event of so-called electrostatic discharges (ESD). High voltages and high currents, associated with transient or ESD disturbances, cause high disturbing powers to occur.\nSuch a pulse (e.g. burst) may also occur in many fields of application, e.g. automotive engineering. In automotive engineering, by way of example, there is the requirement that circuits of this type which have to function in the high-voltage range up to 90 volts or above are also to be designed for significantly higher disturbance pulse levels.\nFor high-voltage applications produced by means of high-voltage processes, provision is usually made of protection devices that are initiated or triggered by an electrical breakdown. The breakdown voltages must be significantly greater than the maximum permissible operating voltages of the application circuit to be protected. Only then is it possible to guarantee an undisturbed functionality of the integrated circuit. In the case of a fault, e.g. when an impermissibly high voltage is present, this overvoltage is dissipated to reference potential or ground by the protection circuit and downstream assemblies are thus protected from the high voltage.\nOne alternative for such breakdown-based protection concepts is an active circuit for protecting an integrated circuit, comprising a combination of an active trigger circuit with a known protection device such as a thyristor or a bipolar or MOS protection transistor.\nActive circuits for protecting the IC are often triggered by the rise of the transient signal. In this case, the signal rise per unit time is detected and a protection transistor or a protection circuit is turned on by means of a drive circuit.\nIn the case of a fault, the protection circuit may accordingly be understood as an actively triggered overvoltage or overcurrent surge arrestor. A fast activation of the protection circuit is necessary in the case of a fault.\nShort switch-on times and a precise switch-on threshold of the protection circuit for the integrated circuit and the protective effect thereof for different forms of disturbance pulses are important aspects of the product specification and constitute a competitive advantage.\nU.S. Pat. No. 5,982,601 discloses a thyristor (SCR—silicon control rectifier) for ESD protection, which is triggered directly by the transient signal. The thyristor is realized in the semiconductor arrangement in a manner known per se by means of an n-type well, a p-type well and highly doped n-type and p-type regions. The transient voltage is detected by means of an RC element. By means of inverters connected downstream, the voltage level detected at the capacitance is converted into a control signal which drives the base of the pap transistor of the thyristor structure. As soon as the output current of the then active pnp transistor generates a sufficiently large voltage drop at a resistor, the npn transistor of the thyristor structure turns on, so that the transient pulse is dissipated from the pad potential of the I/O pin to reference potential through the low-impedance thyristor path. The thyristor then automatically remains turned on until its current falls below the holding current and the turn-off condition is met."}
{"text": "The present invention relates to a novel drug delivery system. More particularly, the present invention relates to a product in which a biologically active material is present in a multiphase system, i.e., (a) captured in multilamellar lipid vesicles (MLV); (b) dissolved in the solvent components of the system; and (c) in a solid crystalline or amorphous state.\nLiposomes are lipid vesicles composed of membrane-like lipid layers surrounding aqueous compartments. Liposomes are widely used to encapsulate biologically active materials for a variety of purposes, but particularly they are used as drug carriers. Depending on the number of lipid layers, size, surface charge, lipid composition and methods of preparation various types of liposomes have been utilized.\nMultilamellar lipid vesicles (MLV) were first described by Bangham, et al., (J. Mol. Biol. 13: 238:252, 1965). A wide variety of phospholipids form MLV on hydration. MLV are composed of a number of bimolecular lamellae interspersed with an aqueous medium. The lipids or lipophilic substances are dissolved in an organic solvent. The solvent is removed under vacuum by rotary evaporation. The lipid residue forms a film on the wall of the container. An aqueous solution generally containing electrolytes and/or hydrophilic biologically active materials are added to the film. Agitation produces larger multilamellar vesicles. Small multilamellar vesicles can be prepared by sonication or sequential filtration through filters with decreasing pore size. Small unilamellar vesicles can be prepared by more extensive sonication. An improved method of encapsulating biologically active materials in multilamellar lipid vesicles is described in U.S. Pat. No. 4,485,054.\nUnilamellar vesicles consist of a single spherical lipid bilayer entrapping aqueous solution. According to their size they are referred to as small unilamellar vesicles (SUV) with a diameter of 200 to 500 .ANG.; and large unilamellar vesicles (LUV) with a diameter of 1000 to 10,000 .ANG.. The small lipid vesicles are restricted in terms of the aqueous space for encapsulation, and thus they have a very low encapsulation efficiency for water soluble biologically active components. The large unilamellar vesicles, on the other hand, encapsulate a high percentage of the initial aqueous phase and thus they can have a high encapsulation efficiency. Several techniques to make unilamellar vesicles have been reported. The sonication of an aqueous dispersion of phospholipid results in microvesicles (SUV) consisting of bilayer or phospholipid surrounding an aqueous space (Papahadjopoulos and Miller, Biochem. Biophys. Acta., 135: 624-238, 1968). In another technique (U.S. Pat. No. 4,089,801) a mixture of a lipid, an aqueous solution of the material to be encapsulated, and a liquid which is insoluble in water, is subjected to ultrasonication, whereby liposome precursors (aqueous globules enclosed in a monomolecular lipid layer), are formed. The lipid vesicles are then prepared by combining the first dispersion of liposome precursors with a second aqueous medium containing amphiphilic compounds, and then subjecting the mixture to centrifugation, whereby the globules are forced through the monomolecular lipid layer and forming the bimolecular lipid layer characteristic of liposomes.\nAlternate methods for the preparation of small unilamellar vesicles that avoid the need of sonication are the ethanol injection technique (S. Batzri and E. D. Korn, Biochem. Biophys. Acta. 298: 1015-1019, 1973) and the ether injection technique (D. Deamer and A. D. Bangham, Biochem. Biophys. Acta. 443: 629-634, 1976). In these processes, the organic solution of lipids is rapidly injected into a buffer solution where it spontaneously forms liposomes--of the unilamellar type. The injection method is simple, rapid and gentle. However, it results in a relatively dilute preparation of lipsomes and it provides low encapsulation efficiency. Another technique for making unilamellar vesicles is the so called detergent removal method (H. G. Weder and O. Zumbuehl, in \"Liposome Technology\" ed. G. Gregoriadis, CRC Press Inc. Boca Raton, Fla., Vol. I, Ch. 7, pg 79-107, 1984). In this process the lipids and additives are solubilized with detergents by agitation or sonication yielding defined micelles. The detergents are then removed by dialysis.\nMultilamellar vesicles can be reduced both in size and in number of lamellae by extrusion through a small orifice under pressure, e.g., in a French press. The French press (Y. Barenholz; S. Amselem and D. Lichtenberg, FEBS Lett. 99: 210-214, 1979), extrusion is done at pressures of 20,000 lbs/in at low temperature. This is a simple, reproducible, nondestruction technique with relatively high encapsulation efficiency, however it requires multilamellar liposomes as a starting point, that could be altered to oligo- or unilamellar vesicles. Large unilamellar lipid vesicles (LUV) can be prepared by the reverse phase evaporation technique (U.S. Pat. No. 4,234,871, Papahadjopoulos). This technique consists of forming a water-in-oil emulsion of (a) the lipids in an organic solvent and (b) the substances to be encapsulated in an aqueous buffer solution. Removal of the organic solvent under reduced pressure produces a mixture which can then be converted to the lipid vesicles by agitation or by dispersion in an aqueous media.\nU.S. Pat. No. 4,016,100, Suzuki, et al., describes still another method of entrapping certain biologically active materials in unilamellar lipid vesicles by freezing an aqueous phospholipid dispersion of the biologically active materials and lipids. All the above liposomes, made prior to 1983, can be classified either as multilamellar or unilamallar lipid vesicles. A newer type of liposomes is referred to as multivesicular liposomes (S. Kim, M. S. Turker, E. Y. Chi, S. Sela and G. M. Martin, Biochim. Biophys. Acta 728; 339-348, 1983). The multivesicular liposomes are spherical in shape and contain internal granular structures. A lipid bilayer forms the outermost membrane and the internal space is divided up into small compartments by bilayer septrum. This type of liposomes required the following composition: an amphiphatic lipid with net neutral charge, one with negative charge, cholesterol and a triacylglycerol. The aqueous phase containing the material to be encapsulated is added to the lipid phase which is dissolved in chloroform and diethyl ether, and a lipid-in-water emulsion is prepared as the first step in preparing multivesicular liposomes. Then a sucrose solution is shaken with the water-in-lipid emulsion; when the organic solvents are evaporated liposomes with multiple compartments are formed.\nFor a comprehensive review of types of liposome and methods for preparing them refer to a recent publication \"Liposome Technology\" Ed. by G. Gregoriadis. CRC Press Inc., Boca Raton, Fla., Vol. I, II, & III 1984.\nSolutions are one of the oldest type of pharmaceutical dosage forms or drug delivery systems. A true solution is defined as a mixture of two or more components that form a homogeneous molecular dispersion, i.e., a one phase system. According to the United States Pharmacopeia, Twentieth Revision (USP XX, page 1027), solutions are liquid preparations that contain one or more soluble chemical substances usually dissolved in water. Further, solutions are used for the specific therapeutic effect of the solute, either internally or externally. Id.\nSuspensions are preparations of finely divided, undissolved drugs dispersed in liquid vehicles (USP XX page 1030). In this sense a suspension is a heterogenous, two-phase system. Suspensions have been used as drug delivery systems for centuries, for providing an insoluble bioactive ingredient for oral, parenteral and for the topical route of administration. In the present multicomponent, multiphase liposomal system the biologically active substance is present in the solid form dispersed in the aqueous medium both inside and outside the lipid vesicles.\nHydrogels can be any one of a wide variety of synthetic and natural hydrophilic polymers. They are used in pharmaceutical dosage formulation for various purposes, i.e., as viscosity inducing agents for suspensions and ophthalmic solutions; as protective colloids to stabilize emulsions and suspensions; as vehicles for topically applied dosage forms; as controlled-release drug delivery systems (J. D. Andrade (ed) \"Hydrogels\" for Medical and Related Applications: ACS Symposium, Series Nol. 31 ACS Washington, D.C., 1976). A gel is generally a semisolid system of at least two constituents, consisting of a condensed mass enclosing and interpenetrated by a liquid. The gel mass may consist of flocculles of small particles or macromolecules existing as twisted, intermingled, matted strands. The polymer units are often bound together by electrostatic, hydrogen and van der Waal forces. Gels containing water are called hydrogels, those containing organic liquid are called organogels.\nThe hydrophilic polymers in aqueous media exhibit \"pseudoplastic flow\" due to the effect of intermolecular entanglements and the binding of water molecules. When the long randomly coiled polymer chain moves, their solvation layer (water of hydration) are dragged along which increases the resistance to flow, or the viscosity of the solution. This property is often utilized in pharmaceutical formulation to increase the viscosity of the preparation. Recently the hydrogels have been disclosed as useful for a controlled drug delivery system (S. W. Kim, Pharmacy International 4: 90-91 (1983))."}
{"text": "In a conveying system for conveying small articles such as cassettes or pallets, it is frequently necessary to remove the articles from the conveyor for processing. Traditionally, an operator has manually gathered the articles as they pass along the conveyor, and manually stacked the articles in preparation for transporting to the processing equipment. The manual gathering and stacking of the articles requires considerable labor and is a substantial cost factor.\nAutomatic stacking mechanisms have been proposed in the past, such as shown in U.S. Pat. No. 3,470,996. In the stacking mechanism, as shown in the aforementioned patent, articles on a conveyor are pushed laterally by a pusher unit from a conveyor to a position beneath a stacking chute. A lift cylinder then operates to elevate the article upwardly into the chute and the elevated article is then held in the chute by a second plunger which is located opposite the pusher assembly. While devices of the type shown in the aforementioned patent can serve to stack articles one on top of each other, the mechanism requires that the articles slide against each other during movement from the conveyor to the stacking chute and this can cause substantial damage and scratching to the articles, if they are of a delicate nature."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a method and a device for operating a hybrid drive system having multiple drives, and relates particularly to the execution of learning functions in order to adapt regulation of the drives of the hybrid drive system to changed physical behavior of the drives.\n2. Description of the Related Art\nFor previous drives having internal combustion engines, the internal combustion engines are activated by an appropriate engine control unit in order to provide the desired output variables such as rotational speed or torque, for example. Since the physical behavior of the internal combustion engine may change over time due to various influences, for example aging, temperature fluctuations, or the like, adaptation functions are usually implemented in the engine control unit. These adaptation functions may be triggered as a function of operating data, for example the operating time of the internal combustion engine, the distance traveled, the number of driving cycles, or the number of revolutions of the engine shaft, and as a function of a specified operating range of the internal combustion engine. If the specified operating range is operated for a minimum period of time which is necessary for completely executing the adaptation function, specific parameters may be adjusted or adapted with the aid of the adaptation function.\nIn hybrid drive systems, two or more than two types of drives are generally used in parallel. A higher-level hybrid control unit is used to select a subdivision of the torque to be provided, according to a drive strategy, as a function of the driver's specifications and as a function of the instantaneous operating states of the vehicle or of the drive system. Depending on the drive strategy provided, the operating control unit may provide that in the specific operating range of the motor vehicle and in specific ranges of driver specifications the load is distributed between the drives so that the drives are each operated in preferred operating ranges. However, this may result in certain operating ranges no longer being included in a drive, since the drive strategy which is implemented in the hybrid control unit keeps one or more of the drives in a preferred operating range.\nIn the example of the internal combustion engine as the drive, current drive strategies of the hybrid control units provide that the hybrid control units are operated in a medium load range so that, for example, operating modes such as coasting mode or idle speed occur rarely or not at all.\nIf this is the case, the engine control unit is no longer able to carry out the adaptation function, even when a suitable testing time is established, since the internal combustion engine no longer reaches the operating range which is necessary for executing the adaptation function.\nAn object of the present invention, therefore, is to provide a method and a device for operating a hybrid drive system via which, also for a hybrid drive system, it may be ensured that necessary adaptation functions for the individual drives are executed."}
{"text": "1. Field of the Invention\nThe present invention relates to a camera which can record information on a magnetic recording medium formed on a silver salt film (hereinafter referred to as a film), film means loadable to such a camera and a camera system.\n2. Related Background Art\nIn an 8 mm talkie, a magnetic track which is a magnetic recording medium is formed on a film so that sound information is recorded as magnetic record information. In a conventional camera, date and time information of photographing is recorded on a film as optical record information by means of optical imaging.\nRecently, a camera which records information on a magnetic track formed on a 35 mm film to allow recording of information other than optical record information has been proposed.\nHowever, because of contact between a magnetic head and a magnetic tape, powders of magnetic material of the tape deposit on a surface of the magnetic head after a certain period of use and this deteriorates a recording characteristic."}
{"text": "1. Field of the Invention\nThe present invention relates to a keycap, and in particular to a light-guiding method and a light-guiding assembly of a backlight source within the keycap.\n2. Description of Prior Art\nBack light source can be widely used in a liquid crystal display, a signboard or an advertisement board. In addition, it can be also applied to electronic products, such as mobile phone or personal digital assistant. Since the back light module is mounted within the keycap, the user can still smoothly operate the keycaps even in the night or without the illumination of a light source.\nThe conventional keycap is provided with a light-guiding plate in the interior of the key pad assembly, and the periphery or the bottom of the light-guiding plate is provided with a light-emitting diode. When the electronic product is in use, the light-emitting diode is lighted up, and the light is introduced into the light-guiding plate. At this time, with the light-guiding plate, the light is directed to the keycap above the light-guiding plate, so that the surface of the keycap has a light-transmitting effect. In this way, the user can clearly see the character, numeral or symbol on the surface of each keycap.\nRecently, with the development of multi-color light-emitting diodes, they are also widely used in the keycaps of electronic products. The multi-color light-emitting diode generates the lights with various colors under the control of the circuit of the electronic product to, so that the surface of the keycap or some keycap groups can represent a colored back light effect with different colors. Since there is only one layer of light-guiding plate, if several multi-color light-emitting diodes are lighted up at the same time, and the lights with various colors are simultaneously introduced into the light-guiding plate, these lights may be mutually interfered with one another or mixed easily. As a result, the light displayed on the surface of the keycap is not the originally preset light. Therefore, in order to solve the mutual interference or mixing of the lights, adhesion of black patches or coating of inks is applied to the light-guiding plate so as to shield the light. After several multi-color lights are introduced into the light-guiding plate, the surface of the keycap or some keycap groups may represent a backlight effect with various colors and do not generate the interference or mixing of the lights. However, if the adhesion or coating process is not perfectly performed, such light-shielding measure may cause the leakage of light and in turn the waste of time and process in manufacturing the light-guiding plate."}
{"text": "As a kind of integrated circuit (IC) memory devices, flash memories have the function of storing information in an electrically erasable and rewritable manner. For this reason, flash memories are widely used in electronic products including portable computers, cell phones and digital music players. A flash memory has many memory cells that are required to be so arranged as to be adapted to its operations. Each of the memory cells stores a single bit of data.\nFIG. 1 schematically illustrates a memory cell of a conventional flash memory device, and FIG. 2 shows an array of memory cells of the conventional device. The flash memory device includes a plurality of the arrayed memory cells, as well as a plurality of bit lines (BL0, BL1, BL2, BL3, . . . , and BLm), a plurality of word lines (WL) and a plurality of control gates (e.g., CG0, CG1 and so forth), configured to select and provide drive signals to the individual memory cells. As shown in FIG. 1, each of the memory cells includes two memory bits, i.e., bits A and B, and a word line common to the two bits. Each of the memory bits includes a bit line and a control gate. While only two control gates are depicted in FIG. 2, i.e., CG0 and CG1, in practical, the memory may have more control gates (i.e., CG0, CG1, CG2, CG3, . . . , and CGn), as generally and commonly known in this art, and a detailed description thereof is therefore omitted herein.\nConventionally, in order to check whether a memory bit has been properly programmed, programming verification is generally carried out by reading the memory bit. This process is explained below with the programming of the memory bit A as an example. Typically, the memory bit A is programmed, concurrently with voltages of the control gate CG0, the word line WL and the control gate CG1 at 8.6 V, 1.5 V and 5V, respectively. In addition, after completion of the programming, the memory bit A is typically verified with the voltages of the control gate COO, the word line WL and the control gate CG1 at 0 V, 4.5 V and 5 V, respectively.\nReferring to FIG. 3, in which the abscissa axis represents time t. The memory bit A is programmed at an instant t0, with the voltage of the control gate CG0 at 8.6 V and the voltage of the word line WL at 1.5 V. After the programming is completed at an instant t1, in order to enable the programming verification, it is needed to decrease the voltage on the control gate CG0 to 0 V and increase the voltage on the word line WL to 4.5V. However, due to latch-up and other effects, the decrease of the voltage at the control gate CG0 to 0 V takes a period of time tf and ends at an instant t2, which lead to elongation of the programming verification time. In addition, after the programming verification has ended at an instant t3, the voltage on the control gate CG0 is required to increase again to 8.6 V, with the voltage on the word line WL required to drop to 1.5 V, in order to allow the next programming cycle to start. Again, due to latch-up and other effects, the increase of the voltage on the control gate CG0 takes a period of time tr and ends at an instant t4. This makes an additional contribution to the elongation of the needed programming time."}
{"text": "A technique in which a transistor is formed using a semiconductor layer formed over a substrate having an insulating surface is known. For example, a technique in which a transistor is formed over a glass substrate using a thin film including a silicon-based semiconductor material and the transistor is applied to a liquid crystal display device or the like is known.\nA transistor which is used in a liquid crystal display device is generally formed using a semiconductor material such as amorphous silicon or polycrystalline silicon. Although transistors including amorphous silicon have low field effect mobility, a larger glass substrate can be used in the case of using amorphous silicon. Meanwhile, transistors including polycrystalline silicon have high field effect mobility; however, they need to be subjected to a crystallization step such as laser annealing and thus are not always suitable for larger glass substrates.\nAs another material, an oxide semiconductor attracts attention. Zinc oxide and a materials including zinc oxide are known as oxide semiconductor materials. Thin film transistors each of which is formed using an amorphous oxide (an oxide semiconductor) having an electron carrier concentration of lower than 1018/cm3 are disclosed in Patent Documents 1 to 3."}
{"text": "TECHNICAL FIELD OF THE INVENTION\nThe present invention relates to novel amide-derived biaromatic compounds and to pharmaceutical/cosmetic compositions comprised thereof; the subject amido compounds are especially useful in human or veterinary medicine, or alternatively in cosmetic compositions."}
{"text": "1. Field of the Invention\nThe present invention relates to verifying illegitimate non-human users that are accessing content.\n2. Description of the Related Art\nIn the past few years, there has been a significant increase in the number of automated non-human user software, known as “bots”, browsing the internet. Some of these bots are used for legitimate purposes to analyze and classify content across the World Wide Web. For example, GOOGLE® uses bots to gather content to be indexed for their search services. However, some other types of bots are used for illegitimate and many times for fraudulent purposes. One such illegitimate usage is the artificial inflation of impression counts (number of times an advertisement is viewed) and/or impression clicks (number of times an advertisement is clicked) in order to fraudulently profit from getting paid based on those inflated numbers.\nThese bots are very difficult to identify because they may originate from a server farm or from regular user computers, computers that real and unsuspecting humans use to legitimately view web pages or other types of digital content. The bots can spread and infect a computer through malware, adware, malvertising, viruses, plugins, email attachments, apps, websites, or through any other means. Further, once potential bots are identified, it is important to verify that that the potential bot actually is a bot and hasn't been misidentified.\nA need arises for effective techniques that can be used to verify illegitimate non-human users that are accessing content."}
{"text": "1. Field of the Invention\nThe present invention concerns a coil for a magnetic resonance imaging system, and particularly an arrangement for fixing such a coil in place.\n2. Description of the Prior Art\nIn a magnetic resonance system, there is a need to use a strong electromagnetic field created by high-powered radio frequency signals so as to excite nuclei in a human body to generate magnetic resonance signals, and then the magnetic resonance signals emitted from the human body are received. The said magnetic resonance signals are reconstructed into magnetic resonance images in a known manner. When a body is examined, normally it is necessary to position one or more radio frequency coils on a patient's bed with their field of view oriented at different parts of the human body, and a suitable coil needs to be placed at a corresponding position in order to carry out a scanning.\nFIGS. 1 to 3 show schematic diagrams of the structure of a currently available breast coil and the manner of fixing it onto a patient bed. Tabs 14 are arranged respectively on both sides 12 of a breast coil 10 in a direction facing a patient bed 20. A plurality of slots 22 are formed respectively on both sides of the patient bed 20 along its longitudinal direction. When the breast coil 10 is used to carry out an examination, said tabs 14 are inserted along the arrow direction into the corresponding slots 22 according to the position that the patient is in said patient bed 20 or the patient's body height, so that the breast coil 10 can be fixed onto said patient bed 20.\nFIGS. 4 to 6 show schematic diagrams of the structure of a currently available head coil and the manner of fixing it onto a patient bed. A number of studs 34 are arranged respectively on the bottom surface 32 of a head coil 30 in a direction facing a patient bed 20. A number of corresponding holes 24 are formed on the patient bed 20. When the head coil 30 is used to carry out an examination, the studs 34 are inserted along the arrow direction into the corresponding holes 24, so that the head coil 30 can be fixed onto the patient bed 20.\nIt can be seen that in such a currently available magnetic resonance imaging system, in order to fix different coils at different positions on the patient bed, it is necessary to arrange correspondingly specific structures, such as tabs and slots, studs and holes, etc., on the various coils and on the patient bed, so as to carry out the positioning or fixing therebetween.\nSince specific structures have to be designed on a patient bed, such as forming slots, holes, etc., the strength of the bed base may be affected. In addition, the structures such as slots, holes, etc. tend to be soiled and collect dirt easily, therefore it is necessary for them to be cleaned regularly to prevent disease transmission.\nIn order to use a coil for scanning at a variety of positions, it is necessary to make a series of slots or holes on the patient bed, and this further reduces the structural strength of the patient bed. Furthermore, such a coil can only be located on the positions where the slots and holes are formed, and this leads to inconvenient placement operations, and also restricts the flexibility in using the coil."}
{"text": "The present invention relates generally to propellant charges in medium and large caliber cased or caseless gun systems. More specifically the invention is directed to the configuration of the charge that is loaded into the cartridge case assembly. The configuration, in this case, generally the shape, density and placement of the propellant charge is a departure from contemporary configurations.\nMunitions that use solid propellants carried in the cartridge casing typically relied on spherical or cylindrical grain geometry for ease of processing. Burn rates, the speed that a flame front is propagated through the propellant, were controlled through geometry (i.e. grain perforations) or through chemical (i.e. deterrents) means. Burn rate modification and progressivity tuning are important considerations in designing charge loads. Another technique to provide a propellant in the cartridge is to use what are known as co-layered propellants in a disk or rolled scroll. Co-layered propellants being a descriptive term describing the use of layers of precut propellant elements cut from a sheet of propellant material. It is known to use co-layered propellants in the shape of rings or appertured disks of different propellants stacked in a sandwich configuration. In co-layered propellants so called \"fast burn\" propellant sheet stock, cut in the appropriate appertured disk or ring shape, is sandwiched between two layers of slower burning propellant sheet stock. The source of co-layered propellant sheets is Alliant Tech Systems who can provide the co-layered propellant in die-cut annular ring shapes or in uncut sheets.\nAnother use of co-layered propellants is in a rolled scroll form. In this configuration the propellant sheet is rolled to form a scroll. The scroll form is aligned longitudially with the boom of the projectile and if ignited from the end of the scroll, the normal technique, the flame propagation will be more difficult to control than the flame propagation of this invention. In the scroll form of propellant it is difficult to control burn propagation and there is a risk that waves of pressure may build to the point of weapon destruction. Therefore the scroll form of propellant is more difficult to specify than is the strip embodiment presented here.\nOne of the disadvantages of co-layered propellants is that the stacked disks provide, having only limited natural ullage, a barrier to gas flow which can lead to detrimental pressure waves in the gun. These detrimental waves are caused in part by the blockage presented by the disks. The invention presented in this specification avoids or minimizes the barrier to the gas flow as the strips of propellant allow sufficient ullage between the strips.\nAnother disadvantage of co-layered propellants in disk form, as is used in the prior designs, is that the manufacture of the rings creates a large amount of scrap propellant. It is known that to cut a circular disc from a rectangular sheet of propellant material will cause significant selvage or waste of material. The sheets of propellant are expensive and anytime you can save on material the cost of a loaded cartridge case will be commensurately lower."}
{"text": "The invention relates to a work-holding fixture for clamping a workpiece in position on a work-holding table, comprising a guide shoe, which can be fixed on the work-holding table and has recesses in symmetric arrangement, a clamping dog, which is designed as a rocking lever, has a clamping leg and a gripping leg and, with laterally projecting extensions arranged in the center region, engages into the recesses of the guide shoe and is supported there during clamping, and a clamping screw which passes through the clamping leg in the region of its free end and is supported on the guide shoe during clamping. Such work-holding fixtures are used in particular in the machining of workpieces by the workpiece being clamped in position on the work-holding table.\nFrench Preliminary Published Specification 2,612,439 discloses a work-holding fixture of the type described at the beginning. In this work-holding fixture, a clamping dog is already allocated to a guide shoe. The guide shoe has two oblong-hole-like apertures and is clamped in position by means of two fixing screws in tapped holes which are provided in the work-holding table, the oblong-hole-like apertures enabling an infinitely variable displacement of the guide shoe in a single horizontal direction. The effective direction of the clamping dog in the horizontal direction is thus also established. The clamping dog is of plate-shaped design and, roughly in its center area, has two laterally projecting extensions, the plate-shaped clamping dog thus being designed as a rocking lever and having a clamping leg extending to the rear and a gripping leg extending forward from the extensions, with which gripping leg the workpiece is ultimately clamped in position. The laterally projecting extensions have a cylindrical cross-section and engage on their part into recesses of the guide shoe. These recesses, with the exception of an insertion shaft which can be closed by a headless setscrew or the like, are designed to be enclosed at the margins and enable the clamping dog to be displaced a certain limited distance in the single direction described above relative to the fixed guide shoe before the workpiece is clamped in position. The adjustability in this direction is only comparatively small. The recesses have an effective length which is arranged within the limits of the oblong-hole-like apertures and is smaller than the length of the oblong-hole-like apertures. A clamping screw passes through the free end of the gripping leg, which clamping screw is guided in a corresponding thread in the gripping leg and, when the workpiece is clamped in position, presses down on the rear end of the guide shoe. The clamping range of this known work-holding fixture is restricted vertically to a few bridgeable millimeters, since the clamping dog is of plate-shaped design and is pivotably mounted about the journal-like extensions essentially only to a very limited extent. Furthermore, it is disadvantageous that, due to the fixing of the guide shoe on the work-holding table, only a single direction in the horizontal plane is predetermined in which the work-holding fixture is infinitely adjustable. The displacement travel is here very small and it is necessary to provide a comparatively close grid of tapped holes in the work-holding table so that workpieces protruding to various widths can be clamped at the various locations in this single direction with one and the same work-holding fixture. In this arrangement, it is also disadvantageous that, as a result of the close grid, the guide shoe often has to be shifted on the work-holding table if different sizes of workpieces require this. Clamping in a horizontal direction other than that predetermined by the two fixing screws is not possible. The guide shoe cannot be rotated about a vertical axis. Furthermore, unfavorable force-transmission ratios also result if, for example, the fixing screws are inserted relatively far to the rear at the oblong-hole-like apertures and the extensions and thus the clamping dog come into use in a position displaced relatively far forward in the recesses. Furthermore, it is disadvantageous that the clamping screw, with the free end of its threaded shank, rests in a planar manner on the guide shoe only at a single, fixed workpiece height. At every deviating height of the workpiece, the clamping screw is supported only in a point-like manner on the guide shoe, in the course of which inadmissibly large surface pressures develop. With this known work-holding fixture, it is not possible to clamp workpieces of different height within a larger range, since the tilt angle of the plate-shaped clamping dog is limited to a few degrees.\nU.S. Pat. No. 4,470,586 shows a work-holding fixture in which there is no guide shoe. Instead, the clamping dog, which here, too, has two laterally projecting extensions interacts directly with the slot of a slotted work-holding table. For this purpose, the clamping dog has a bent shape, the gripping leg and the clamping leg extending roughly in a straight line relative to one another and the extensions being arranged well below this straight line, since they have to engage into the slot of the work-holding table. A clamping screw is also provided here at the rear end of the clamping dog, that is, at the free end of the clamping leg, which clamping screw is supported on the root of the slot in the slotted table. The laterally projecting extensions provided on the clamping dog have a slightly oval shape, as a result of which the clamping height is limited. The clamping height ends here when the oval extensions abut against two opposite surfaces of the slot of the slotted table; higher workpieces can no longer be clamped with this work-holding fixture. It is thus possible to alternatively clamp workpieces in a height range which exceeds a few centimeters. In the horizontal plane, however, the work-holding fixture is tied to the direction which is predetermined by the direction of the slots in the slotted work-holding table. The work-holding fixture cannot be used in an angular position turned, for example, through 90.degree. thereto or also in any oblique angular position. The free end of the clamping screw here is certainly of spherical design; but here, too, the disadvantage of the point-like force transmission results. A precondition for the applicability of this work-holding fixture is a work-holding table provided with slots. In machining centers, the work-holding tables provided there are often only provided with tapped holes which are arranged in a certain grid. The known clamping dog cannot be used in combination therewith."}
{"text": "Liquid crystal displays have been increasingly becoming popular as displays for information, image and the like. The liquid crystal displays have advantages of providing thin profile, small placing area and the like when compared with CRT (Cathode Ray Tube) displays.\nLiquid crystal displays according to the present invention can be used in, for example, a display or monitor for a personal computer, a monitor for factory automation (FA), a home television set, a terminal monitor in a hospital, a library, a museum and the like, a monitor for an air traffic control tower and the like, a monitor used for reference of newspapers, documents in public offices and the like, a personal monitor in a school and a supplementary private school, a personal monitor for utilizing various media, a monitor in amusement facilities such as pachinko facilities, and the like. The liquid crystal displays according to the present invention can also be used as a light valve for projector type liquid crystal displays.\nA recent liquid crystal display has image quality substantially equal to that of a CRT display. However, a viewing angle of the liquid crystal display is much narrower than that of the CRT display. Therefore, enlarging a viewing angle is one of the most important concerns in a liquid crystal display.\nAs a technology for enlarging a viewing angle of a liquid crystal display, there is considered using a liquid crystal display element which has a plurality of domains having different alignment directions in each pixel to obtain a wide viewing angle. Such prior art is disclosed in, for example, Japanese examined patent publication No. 58-43723, “Liquid Crystal Display Element”, Japanese patent laid-open publication No. 59-211019, “Liquid Crystal Display Device”, and Japanese patent laid-open publication No. 63-106624 (Japanese patent publication No. 2692693), “Liquid Crystal Display Panel”.\nReferring to a “Liquid Crystal Display Panel” disclosed in Japanese patent laid-open publication No. 63-106624 as a prior art example, among the Japanese publications mentioned above, an explanation will be made on a prior art technology.\nFIG. 28 is a plan view showing a conventional liquid crystal display according to this prior art example. FIG. 28 corresponds to FIG. 1 described in the Japanese patent laid-open publication. No. 63-106624. FIG. 29 corresponds to FIG. 2 described in the Japanese patent laid-open publication No. 63-106624. FIG. 30 corresponds to FIG. 3 described in the Japanese patent laid-open publication No. 63-106624. FIG. 30 is a cross sectional view taken on line a-a′ of FIG. 28, and shows a cross section of one pixel of the conventional liquid crystal display device, that is, a liquid crystal display element.\nWith reference to FIGS. 28 through 30, a structure of the conventional liquid crystal display device will be described.\nOn a lower glass substrate 122, there are formed a transparent lower display electrode 120 which is provided corresponding to each pixel, and a lower alignment film or layer 110 formed on the transparent lower display electrode 120.\nOn the lower glass substrate 122, there is also formed a thin film transistor (TFT) 113 for driving the transparent lower display electrode 120.\nOn an upper glass substrate 121, there are formed a transparent upper display electrode 111 and a upper alignment film or layer 109. The alignment films 109 and 110 are both made of polyimide.\nA pixel B formed between the opposing transparent electrodes 111 and 120 has, for example, a square shape of 200 μm (micrometers) square, and a plurality of such pixels B are disposed in a matrix. At the center portion of the transparent display electrode, there is disposed a band shaped spacer 123 made of polyimide. As a result thereof, each pixel B is divided into a region X131 and a region Y132 by the band shaped spacer 123. These divided regions X131 and Y132 are formed as shown schematically in FIG. 16 which corresponds to FIG. 2 of Japanese patent laid-open publication No. 63-106624. That is, the upper glass substrate 121 and the opposing lower glass substrate 122 are respectively rubbed in directions shown by arrows.\nIn order to show an effect of improvement in viewing angle dependence of viewing quality realized by this prior art liquid crystal display, FIG. 31 shows a plan perspective view which illustrates rubbing directions and direction of twist of liquid crystal alignment between the upper and lower substrates when a direction of observation is changed by 45 degrees in bearing angle. In FIG. 31, there are shown divided regions X131, divided regions Y132, rubbing directions 1311 and 1312 in a first substrate, rubbing directions 1312 and 1322, and a twist angle 127 of liquid crystal alignment between the upper and lower substrates.\nFIG. 32 is a cross sectional view taken on line b-b′ of FIG. 31 and showing a rising direction of aligning force at the surface of a substrate and a rising direction of liquid crystal alignment by an electric field between the upper and lower substrates. In FIG. 32, a twist of liquid crystal molecules is not illustrated. In FIG. 32, there are shown liquid crystal material 105, a direction of rising 106 of liquid crystal alignment caused by an electric field, the upper glass substrate 121, a lower glass substrate 122, a region X131, and a region Y132.\nIn this prior art liquid crystal display, as shown in FIG. 31, with respect to liquid crystal alignment in each of the divided regions or domains, a direction of spiral twist is the same, but as shown in FIG. 32, directions of angles from the substrate surface, i.e., pretilt angles, differ from each other. Because of the difference of the angles from the substrate surface, when a voltage is applied between the upper and lower electrodes, direction 106 of rising up of liquid crystal molecules, which is also called direction of tilt, differ from each other as shown in FIG. 31. Therefore, when light is irradiated from a direction which is inclined from the direction perpendicular to the substrate, optical characteristics of respective domains compensate with each other. As a result thereof, when the voltage is applied, viewing angle dependence is compensated with each other between domains having mutually different alignments in each pixel between the upper and lower substrates, and optical characteristics having small viewing angle dependency can be realized. Especially, when an image having gray shades is displayed, a phenomenon of gray shade inversion is not observed.\nOther liquid crystal display panels are also known which function based on the same principle as that of the above-mentioned liquid crystal display panel and which can be fabricated by smaller number of steps than that of the above-mentioned liquid crystal display panel. Such liquid crystal display panels are described in “SID '92 Digest of U.S.A., p. 798”, “Japan Display '92 Digest, p. 591”, “SID '93 Digest of U.S.A., p. 269”, and the like.\nIn these examples, as in the first example, with respect to alignment of molecules of the liquid crystal material in each of the divided domains, direction of twist of a spiral is the same, but an angle to the substrate surface differs from each other. However, these liquid crystal display panels differ from the first example in a direction of aligning force applied to molecules of the liquid crystal material at the surface of the substrate and magnitude of angle of the aligning force.\nIn order to clarify the difference between these examples and the first example, perspective plan views and cross sectional views for each of these examples. The perspective plan views show directions of rubbing and directions of twist of alignment of liquid crystal molecules between upper and lower substrates. The cross sectional views show directions of rising up of the alignment force at the surface of the substrate and directions of rising up of alignment of liquid crystal between the upper and lower substrates caused by an electric field.\nFIG. 33 is a plan perspective view showing directions of rubbing and directions of twist of alignment of liquid crystal molecules between upper and lower substrates, in the prior art example shown in “SID '92 Digest of U.S.A., p. 798”. FIG. 34 is a cross sectional view taken on line c-c′ of FIG. 33, and shows directions of rise of the aligning force at the surface of the substrate and directions of rise of alignment of liquid crystal molecules between the upper and lower substrates caused by an electric field.\nAn explanation will be made on the liquid crystal display element shown in FIG. 33 and FIG. 34. In the plan perspective view of FIG. 33, there are shown a divided domain X131, a divided domain Y132, a rubbing direction 1311 in the domain X on an upper glass substrate, a rubbing direction 1312 in the domain X on a lower glass substrate, a rubbing direction 1321 in the domain Y on the upper glass substrate, a rubbing direction 1322 in the domain Y on the lower glass substrate, and twist angle 127 of the liquid crystal alignment. Also, in the cross sectional view of FIG. 34, there are shown the upper glass substrate 121, the lower glass substrate 122, the divided domain X131, the divided domain Y132, liquid crystal material 105, direction 106a of rise of alignment of liquid crystal molecules caused by an electric field in the domain X, and direction 106b of rise of alignment of liquid crystal molecules caused by an electric field in the domain Y.\nIn this example, rubbing direction of the upper substrate is one direction, i.e., the same direction, in both of the divided domains, and rubbing direction of the lower substrate is one direction, i.e., the same direction, in both of the divided domains. However, materials of the upper and lower substrates are changed between the domain X and the domain Y. Therefore, as shown in FIG. 34, in the domain X131, the angle of liquid crystal molecules from the surface of the lower substrate 122 is higher (larger) than the angle of liquid crystal molecules from the surface of the upper substrate 121. In the domain Y132, the angle of liquid crystal molecules from the surface of the lower substrate 122 is lower (smaller) than the angle of liquid crystal molecules from the surface of the upper substrate 121. Therefore, in the domain X131, the direction 106a of rise of liquid crystal alignment caused by an electric field becomes a direction of aligning force at the lower glass substrate 122, when a voltage is applied to this liquid crystal display element. In this case, in the domain Y132, the direction 106b of rise of liquid crystal alignment caused by an electric field becomes a direction of aligning force at the upper glass substrate 121.\nFIG. 35 is a plan perspective view showing directions of rubbing and directions of twist of alignment of liquid crystal molecules between upper and lower substrates, in the prior art example shown in “Japan Display '92 Digest, p. 591”. FIG. 36 is a cross sectional view taken on line d-d′ of FIG. 35, and shows directions of rise of the aligning force at the surface of the substrate and directions of rise of alignment of liquid crystal molecules between the upper and lower substrates caused by an electric field.\nAn explanation will be made on the liquid crystal display element shown in FIG. 35 and FIG. 36. In the plan perspective view of FIG. 35, there are shown a divided domain X131, a divided domain Y132, a rubbing direction 1311 in the domain X on an upper glass substrate, a rubbing direction 1312 in the domain X on a lower glass substrate, a rubbing direction 1321 in the domain Y on the upper glass substrate, a rubbing direction 1322 in the domain Y on the lower glass substrate, and twist angle 127 of the liquid crystal alignment. Also, in the cross sectional view of FIG. 36, there are shown the upper glass substrate 121, the lower glass substrate 122, the divided domain X131, the divided domain Y132, liquid crystal material 105, direction 106 of rise of alignment of liquid crystal molecules caused by an electric field.\nIn this example, rubbing direction of the upper substrate is one direction, i.e., the same direction, in both of the divided domains, and rubbing directions of the lower substrate have two different directions. Also, as shown in FIG. 36, in both the domain X131 and the domain Y132, the angle of liquid crystal molecules from the surface of the lower substrate 122 is higher (larger) than the angle of liquid crystal molecules from the surface of the upper substrate 121. However, in the lower glass substrate 122, directions of aligning force differ between the domain X131 and the domain Y132. Therefore, in both the domain X131 and the domain Y132, the direction 106 of rise of alignment of liquid crystal molecules caused by an electric field becomes a direction of aligning force at the lower glass substrate 122, when a voltage is applied to this liquid crystal display element.\nFIG. 37 is a plan perspective view showing directions of rubbing and directions of twist of alignment of liquid crystal molecules between upper and lower substrates, in the prior art example shown in “SID '93 Digest of U.S.A, p. 269”. FIG. 38 is a cross sectional view taken on line e-e′ of FIG. 37, and shows directions of rise of the aligning force at the surface of the substrate and directions of rise of alignment of liquid crystal molecules between the upper and lower substrates caused by an electric field.\nAn explanation will be made on the liquid crystal display element shown in FIG. 37 and FIG. 38. In the plan perspective view of FIG. 37, there are shown a divided domain X131, a divided domain Y132, a rubbing direction 1311 in the domain X on an upper glass substrate, a rubbing direction 1312 in the domain X on a lower glass substrate, a rubbing direction 1321 in the domain Y on the upper glass substrate, a rubbing direction 1322 in the domain Y on the lower glass substrate, and twist angle 127 of the liquid crystal alignment. Also, in the cross sectional view of FIG. 38, there are shown the upper glass substrate 121, the lower glass substrate 122, the divided domain X131, the divided domain Y132, liquid crystal material 105, direction 106 of rise of alignment of liquid crystal molecules caused by an electric field.\nIn this example, rubbing direction of the upper substrate is one direction, i.e., the same direction, in both of the divided domains, and rubbing direction of the lower substrate is also one direction, i.e., the same direction, in both of the divided domains. Also, in both the domain X131 and the domain Y132, the angle of liquid crystal molecules from the surface of the lower substrate 122 is the same as the angle of liquid crystal molecules from the surface of the upper substrate 121. However, in each of the domain X131 and the domain Y132, directions of aligning force differ between the upper substrate 121 and the lower substrate 122. When no voltage is applied to this liquid crystal display element, alignment of liquid crystal molecules at the central portion is held horizontal. When a voltage is applied to this liquid crystal display element, the direction 106 of rise of alignment of liquid crystal molecules is determined by electric field having slant directions as shown in FIG. 38.\nIn these three prior art examples, by adopting structural improvement as shown in FIG. 33 through FIG. 38, the number of process steps is decreased from that of the above-mentioned prior art. In any of these methods, as shown in FIG. 34, FIG. 36 and FIG. 38, when a voltage is applied, directions of rise of liquid crystal molecules differ from each other in respective domains. Therefore, when light is irradiated from a direction slanted from the normal of the substrate, optical characteristics of both regions are mutually compensated. As a result thereof, viewing angle dependency in the condition a voltage is applied is mutually compensated between regions having different alignments in each pixel between the upper and lower substrates, so that optical characteristics having reduced viewing angle dependency.\nOn the other hand, there is a commonly known thin film transistor (TFT) type liquid crystal display which uses usual twisted nematic (TN) type liquid crystal display mode. Also, among such commonly known liquid crystal display, there is a liquid crystal display in which liquid crystal alignment is controlled by forming a slit or slits in a pixel electrode of an opposing substrate. Such liquid crystal display is disclosed in Conference Record of the 1991 International Display Research Conference, p. 239 and an explanation will be made on such liquid crystal display.\nFIG. 39, which corresponds to FIG. 1(b) of the above-mentioned document, shows a cross sectional view of a part, i.e., one pixel, of the liquid crystal display disclosed in the above-mentioned document. FIG. 40, which corresponds to FIG. 7(a) of the above-mentioned document, illustrates electric field in lateral direction in such liquid crystal display. Also, FIG. 41, which corresponds to FIG. 1(a) of the above-mentioned document, shows a cross sectional view of a part, i.e., one pixel, of the commonly known liquid crystal display which uses usual TN type liquid crystal display mode. FIG. 29, which corresponds to FIG. 7(b) of the above-mentioned document, illustrates electric field in such liquid crystal display.\nFirst, an explanation will be made on the liquid crystal display shown in FIG. 39 and FIG. 40. In the cross sectional view of FIG. 39, there are shown an upper glass substrate 221, a lower glass substrate 222, a signal electrode line 215, a lower transparent electrode, that is, a pixel electrode, 220, an upper transparent electrode, that is, an opposing electrode 219, a black mask 231, a color layer 232, an overcoat 233, a lower silicon nitride film 204, and an upper silicon nitride film 241.\nOn the upper glass substrate 221, there are disposed, on the side facing toward the lower glass substrate 222, the black mask 231 for shutting out the light, the color layer 232 for passing the light of predetermined color, the overcoat 233 for protecting the black mask 231 and the color layer 232, and the opposing electrode 219.\nOn the lower glass substrate 222, there are disposed, on the side facing toward the upper glass substrate 221, the signal electrode line 215, the lower silicon nitride film 204 for covering the lower glass substrate 222 and the signal electrode line 215, the pixel electrode 220 formed on the lower silicon nitride film 204, and the upper silicon nitride film 241 covering whole portion on the lower glass substrate 222.\nIn FIG. 40, electric field 202 is shown by arrows which is produced when a voltage is applied to the liquid crystal display element having the structure shown in FIG. 39.\nNext, an explanation will be made on the commonly known TN type liquid crystal display element shown in FIG. 41 and FIG. 42. In the cross sectional view of FIG. 41, there are shown an upper glass substrate 221, a lower glass substrate 222, a signal electrode line 215, a lower transparent electrode, that is, a pixel electrode, 220, an upper transparent electrode, that is, an opposing electrode 229, a black mask 231, a color layer 232, an overcoat 233, a lower silicon nitride film 204, and an upper silicon nitride film 241.\nOn the upper glass substrate 221, there are disposed, on the side facing toward the lower glass substrate 222, the black mask 231 for shutting out the light, the color layer 232 for passing the light of predetermined color, the overcoat 233 for protecting the black mask 231 and the color layer 232, and the opposing electrode 229.\nOn the lower glass substrate 222, there are disposed, on the side facing toward the upper glass substrate 221, the signal electrode line 215, the lower silicon nitride film 204 for covering the lower glass substrate 222 and the signal electrode line 215, the pixel electrode 220 formed on the lower silicon nitride film 204, and the upper silicon nitride film 241 covering whole portion on the lower glass substrate 222.\nIn FIG. 42, electric field 202 is shown by arrows which is produced when a voltage is applied to the liquid crystal display element having the structure shown in FIG. 41.\nIn the prior art liquid crystal display element shown in FIG. 39 and FIG. 40, disturbance of liquid crystal alignment, which is called disclination and which is caused by lateral electric field produced between a pixel electrode and a wiring electrode such as a scanning electrode line, a signal electrode line and the like, can be reduced.\nHere, the first problem in the above-mentioned prior art liquid crystal display element is that disturbance of alignment which is called disclination and which is caused by lateral electric field between the wiring electrode and the pixel electrode. This is because, potential of the wiring electrode such as the scanning electrode line, the signal electrode line and the like gives large influence on the electric field. Thereby, relatively large lateral electric field is produced and liquid crystal alignment is changed by such lateral electric field. In order to reduce the disclination, it was necessary to provide a slit in the opposing electrode by patterning it.\nThe second problem in the above-mentioned prior art liquid crystal display element is that, when pressure is applied from the outside to the liquid crystal display panel having such liquid crystal display element by pressing the panel by a finger or by a pen point, alignment of liquid crystal molecules is easily disturbed, or alignment division collapses or becomes defective.\nThis is because, the prior art liquid crystal display panel does not have sufficient mechanical strength and cell gap of the liquid crystal display panel changes easily by the pressure applied thereto, so that alignment of the liquid crystal molecules is changed.\nAlso, it is impossible to easily and simply realize division of a pixel region into a plurality of domains having different alignment to obtain a wide viewing angle.\nThere are other Japanese patent applications each relating to realizing a wide viewing angle in a liquid crystal display.\nOne of such applications is Japanese patent laid-open publication No. 7-104282 and entitled “Liquid Crystal Display and Method of Manufacturing the Same”. This application provides a liquid crystal display and a method of manufacturing the same in which a liquid crystal portion is divided or partitioned by high polymer material into a plurality of liquid crystal regions, and in which liquid crystal molecules or high polymer material in the liquid crystal regions are driven to various modes depending on the alignment condition, or to a mode in which the liquid crystal region is aligned into symmetry with respect to an axis.\nThis liquid crystal display comprises a plurality of liquid crystal regions each of which has an electrode, is disposed between a pair of substrates, at least one of the substrates being transparent, and is surrounded by the high polymer material. This liquid crystal display also comprises a light control layer which is disposed between the surface of at least one substrate and an electrode formed on the substrate, which controls an intensity of incident light by changing transmittance and which has locally different film thickness. Also, the space between the electrodes on respective surfaces of the pair of substrates varies corresponding to the change of the film thickness of the light control layer. By such structure, mechanical strength is improved, and also viewing characteristic, contrast and color reproductivity are improved.\nAnother one of such applications is Japanese patent laid-open publication No. 7-333612 and entitled “Liquid Crystal Display Device”. This application provides a liquid crystal display device having a wide viewing angle by averaging angle dependency of optical transmittance. The liquid crystal display device disclosed in this application has a pair of opposing substrates each of which comprises a transparent electrode formed on the surface thereof and an alignment film formed on the transparent electrode, and a liquid crystal layer of liquid crystal molecules filling the space between the substrates. When a voltage is not applied to the transparent electrodes, in a plane perpendicular to the direction of the thickness of the liquid crystal layer, there are a first liquid crystal domain in which direction of alignment of liquid crystal molecules is approximately constant and a second liquid crystal domain in which direction of alignment of liquid crystal molecules differs from that of the first liquid crystal domain. These two liquid crystal domains are realized by providing uneven portions at alignment layer formed on at least one of the substrate.\nIn a conventional liquid crystal display of a twisted nematic (hereafter referred to as “TN”) type which is widely used, when an electric field is not applied to liquid crystal molecules, a direction of an alignment vector of the liquid crystal molecules is parallel to a surface of a substrate and this condition corresponds to the condition in which “white” is displayed. When a voltage is applied to the liquid crystal display, the direction of the alignment vector varies toward a direction of an electric field according to a voltage applied. Thereby, the condition in which “white” is displayed is gradually changed to the condition in which “black” is displayed.\nHowever, because of behavior peculiar to liquid crystal molecules in a condition a voltage is applied, there occurs a problem that the TN type liquid crystal display has a narrow viewing angle. This problem of narrow viewing angle is prominent in a direction toward which liquid crystal molecules rise when medium gray is displayed.\nAs a method of improving viewing angle characteristics, there are proposed some liquid crystal displays disclosed in Japanese patent laid-open publication No. 4-261522, Japanese patent laid-open publication No. 6-43461 and Japanese patent laid-open publication No. 10-333180.\nIn these liquid crystal displays, viewing angle characteristics are improved by using a structure shown in FIG. 44A, FIG. 44B and FIG. 44C. A liquid crystal display panel of the structure shown in these drawings comprises an upper substrate having a color filter substrate 501 and having a common electrode 502 and an alignment layer 503 which are formed on the color filter substrate and which have a slit 517 formed in the common electrode 502 and the alignment layer 503. Also, there is provided a lower substrate having a substrate 507, a pixel electrode 504 formed on the substrate 507 and an alignment layer 503 formed on the pixel electrode 504. A space between the upper substrate and the lower substrate is filled with liquid crystal molecules 508 which are homeotropic aligned and which form a liquid crystal layer. The liquid crystal display panel shown in FIG. 44A, FIG. 44B and FIG. 44C is put between a pair of polarizers disposed such that polarizing axes or transmission axes of these polarizers have right angles with each other. When a voltage is applied between the common electrode 502 and the pixel electrode 504, an electric field having slant components is produced. Therefore, more than one liquid crystal domains having different alignment directions are produced in each pixel, and thereby viewing angle characteristics are improved. The shape of the slit can be, for example, a rectangle as shown in FIG. 44B, or can be x-shape as shown in FIG. 44C.\nIn the liquid crystal display disclosed in Japanese patent laid-open publication No. 4-261522, directions of inclination of liquid crystal molecules when an electric field is applied thereto are controlled to obtain a high contrast ratio. Also, as disclosed in Japanese patent laid-open publication No. 6-43461, it is possible to use an optical compensation film to improve viewing angle characteristics of black area.\nFurther, in Japanese patent laid-open publication No. 6-43461, not only in a liquid crystal display panel having homeotropic alignment, but also in a liquid crystal display panel having TN alignment, each pixel is divided into more than one domains by using an electric field having slant and different components to improve viewing angle characteristics.\nIn Japanese patent laid-open publication No. 10-333180, a technology is disclosed in which thin film transistors, gate lines and drain lines are disposed under display electrodes, in order to prevent an electric field from the thin film transistors, the gate lines and the drain lines from affecting slant electric field components produced by using a common electrode having opening portions.\nIn Japanese patent laid-open publication No. 10-20323, there is disclosed a liquid crystal display in which more than one kind of minute domains coexist in a liquid crystal layer. One of substrates of liquid crystal display panel of this liquid crystal display has opening portions, and there is provided a second electrode in each of the opening portions. A voltage is applied to the second electrode to produce an electric field having slant components. Thereby, domains having different alignment directions are produced in each pixel and wide viewing angle characteristics are obtained. This publication mainly concerns a liquid crystal display having TN type cell.\nJapanese PCT patent laid-open publication No. 5-505247 discloses a liquid crystal display of In-Plain-Switching (IPS) system or mode in which, in order to rotate liquid crystal molecules while keeping the liquid crystal molecules in a direction parallel to substrates of liquid crystal display panel, two electrodes are disposed on one of the substrates. A voltage is applied to between the two electrodes and, thereby, an electric field parallel to the substrates is produced. In such IPS system, a long axis of each of liquid crystal molecules does not rise from the substrate when the voltage is applied. Therefore, when a viewing angle is changed, variation of birefringence or double refraction of liquid crystal is relatively small and, therefore, wide viewing angle characteristics are obtained.\nAlso, in each of “Journal of Applied Physics”, Vol. 45, No. 12 (1974) 5466 and Japanese patent laid-open publication No. 10-186351, there is disclosed a liquid crystal display which has the above-mentioned IPS mode and in which liquid crystal molecules having a positive anisotropy of permittivity are aligned as homeotropic alignment. The liquid crystal molecules are rotated toward a direction parallel to the substrate by using an electric field parallel to the substrate. In this case, because of directions of the electric field, the liquid crystal molecules which are homeotropic aligned are grouped into two domains. The two domains have different alignment directions of liquid crystal molecules and, therefore, it is possible to obtain a liquid crystal display having wide viewing angle characteristics.\nJapanese patent laid-open publication No. 10-186330 proposes a liquid crystal display in which square wall structures are produced by using a photosensitive material, and a pixel is formed by using each of such structures as a basic unit. When a voltage is applied to produce an electric field, liquid crystal molecules having a negative anisotropy of permittivity are rotated toward different directions in a pixel.\nHowever, in the above-mentioned technology in which slits are provided in the common electrode, it is necessary to perform a microfabrication process, such as photolithography and the like, of the common electrode 502 which is not required in the fabrication process of the usual TN type liquid crystal display. In addition to this, there is a problem that a very precise assemble technology is required for precisely combine the upper substrate 501 and the lower substrate 507. This problem is especially important in an active matrix type liquid crystal display which uses switching elements such as TFT's and the like.\nThat is, in a usual active matrix type liquid crystal display, active elements such as thin film diodes and the like are formed on one of transparent substrates. Therefore, it is only necessary to perform a microfabrication process, such as photolithography and the like, on one of the substrate on which active elements are formed. It is not necessary to perform such microfabrication process on the other substrate which is usually called a “common electrode”, but it is only necessary to form an electrode on whole area of the other substrate.\nHowever, in the above-mentioned prior art technologies, it is also necessary to perform a microfabrication process, such as photolithography and the like, on the “common electrode” on which, usually, it is not necessary to perform such microfabrication process. Therefore, fabrication process becomes complicated, and a high precision alignment and assemble technology is required for the upper substrate 501 and the lower substrate 507.\nWhen the thin film transistors, the gate lines and the drain lines are disposed under the display electrodes, as described in Japanese patent laid-open publication No. 10-333180, an aperture ratio becomes deteriorated.\nAlso, in the liquid crystal display described in Japanese patent laid-open publication No. 10-20323, it is necessary to apply a voltage to each of second electrodes when driving the liquid crystal display and, therefore, a special drive technology is required. Further, an additional process is required in which a voltage is applied to the second electrodes in order to obtain divided alignment directions.\nIn the IPS type liquid crystal display and in the liquid crystal display in which vertically aligned liquid crystal molecules are rotated down by using the lateral electric field, there is a problem that an aperture ratio is deteriorated. Also, in these liquid crystal displays, when a cell gap of each pixel is narrowed to obtain high operating speed, a high drive voltage is required.\nIn the IPS type liquid crystal display and in the liquid crystal display in which vertically aligned liquid crystal molecules are rotated down by using the lateral electric field, there is still another problem as follows. That is, in such prior art liquid crystal displays, since a color filter layer is disposed between a layer in which liquid crystal molecules are disposed and an opposing substrate, especially when switching elements are formed by using a TFT structure, an electric field produced by applying a voltage between the source electrode and the common electrode affects the color filter layer, and deteriorates display characteristics.\nThat is, in coloring matter constituting the color filter layer, sodium ions are included as impurities. Therefore, when an electric field is applied to the color filter, electric charges are stored in the color filter layer and the color filter layer is charged up. If the color filter layer is charged up, an unnecessary electric field is always applied to the liquid crystal molecules under the charged up portion of the color filter layer. Therefore, such electric field affects display characteristics, especially color irregularity.\nIn the liquid crystal display which has square walls disclosed in Japanese patent laid-open publication No. 10-186330, it is necessary to fabricate the walls by using photolithography, in order to realize divided alignment directions. Therefore, a manufacturing process becomes complicated."}
{"text": "Many vessel-mountable sonar systems are known in the art. Examples of such systems include depth sounder(s), simple obstacle avoidance sonar(s), side scan sonar(s), mechanical scanning sonar(s), and manual pole mount(s). The display methods of these vessel-mountable sonar systems do not adequately facilitate navigation with the sonar imageries, and these systems are difficult to understand and challenging to operate while the vessel is moving. In addition, these vessel-mountable sonar systems do not have the capability to provide high definition resolution imaging from a moving platform on the vessel.\nU.S. Pat. No. 7,035,166 discloses a sonar system comprising at least one sonar emitter, an array of sonar sensors configured to receive an echo from an outgoing signal generated by the sonar emitter and generate a raw data signal; at least one roll, tilt, and yaw sensor, an input from a sensor determining the latitude and longitude position of the sonar emitter; and circuitry adapted to collect the raw signal data, process the raw signal data to provide image data. The image is rotated responsively to the roll, tilt and yaw orientation of the sonar system as determined by the respective sensors. The roll, tilt and yaw sensor and image processing allows successive display of multiple images from successive transmissions to be oriented with the same fixed frame of reference in a direction relative to the earth.\nU.S. Pat. No. 7,173,879 discloses a sonar system comprising a broad beam transmitter, a phased array receiver, a processor which utilizes a non-Fourier based beamformer, and correlation based processing techniques to extract valid targets from the beamformed data where the number of targets to search for is based on the values attained from an Adaptive Target Population Estimator processing technique. The values of the eigenvalues and the slope of successive eigenvalues are used to determine the number of targets to be searched for by the beamformer.\nU.S. Pat. No. 7,123,546 discloses a sonar system comprising at least 2 transmitters with varying field of views, a receive system with a field of view which encompasses all the field of views generated by the transmitters, a control system designed to switch electronics between transmitters, and processing parameters set to correspond to the transmitters where the controller allows for the transmission and reception of a single ping before the next transmission, allowing each transmission to be differentiated in time."}
{"text": "The present invention relates to the filtration field, and more particularly, to an improved centrifugal filtration apparatus for filtering and concentrating a solution. The improvement comprises eliminating any retentate pockets below the filter and instead employing a port on the filtrate side of the filter that is above the bottom level of the filter, thereby not permitting filtering to dryness.\nThe filtration of fluids may be accomplished through the use of filtration devices which utilize microporous filters to filter and concentrate a macro-molecular solution is well known. This technique has been used in centrifugal filtration apparatuses that rely on centrifugal forces to create pressure in the apparatus to force solutions through a filter which separates liquid solutions into filtrate and concentrate.\nThere are certain drawbacks, however, associated with conventional centrifugal filtration apparatus. While many apparatus designs can prevent filtering to dryness in certain applications, they do not in other applications, such as a swinging bucket centrifuge. Those that work in both a fixed angle centrifuge and a swinging bucket centrifuge without filtering to dryness may present the problem of either retaining a different amount of concentrate when used in a fixed angle centrifuge than they will when used in a swinging bucket centrifuge or retaining a different amount of concentrate when used in fixed angle centrifuges with different fixed angles.\nSome of the conventional centrifugal filtration apparatus require a second spin to remove the retained concentrate (i.e. dead stop volume). In these apparatus it is difficult or impossible to remove the retained concentrate from the concentrate tube with a pipette.\nAnother problem with some of the conventional centrifugal filtration apparatus is that as the concentrate volume approaches its final retained volume, the active filter area approaches zero. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume.\nAnother problem with the conventional centrifugal filtration apparatus is that they are open systems (i.e. they contain a venting means that vents to the atmosphere in the centrifuge).\nYet another problem with some of the conventional centrifugal filtration apparatus is that they are not scaleable (i.e. they are designed to be used as small volume centrifugal filters, or large volume centrifugal filters, not both).\nCertain types of filtration devices, such as that disclosed in U.S. Pat. No. 4,632,761 to Bowers et alia, are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. This device, however, filters to dryness when spun at a 90xc2x0 angle and therefore the dead stop feature will not work if the device is spun in a swinging bucket centrifuge. Also, in this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. The filter surface area in this type of device is limited by the diameter of the device, and the surface area is relatively small when compared to the volume of liquid solution within the housing. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it contains a venting means that vents to the atmosphere in the centrifuge). This type of filtration device is conductive to clogging because the heaviest and denser molecules within the liquid solution are forced into the membrane filter. Accordingly, this device is limited because it will filter to dryness when spun in a swinging bucket centrifuge, it will filter to different dead stop volumes when used in fixed angle rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, and its active filter surface area is limited by the diameter of the device.\nThe types of filtration devices disclosed in U.S. Pat. No. 4,722,792 to Miyagi et alia are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. In this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. Therefore, the dead stop volume will be different when the device is used in a swinging bucket rotor centrifuge than it will be when used in a fixed angle rotor, and will also be different when the device is used in fixed angle rotor centrifuges of different rotor angle. This means that the results obtained from this type of device when used in one type of centrifuge rotor cannot be compared to the results obtained from this type of device when used in another type of centrifuge rotor. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it must contain a venting means that vents to the atmosphere in the centrifuge). Another problem with this type of centrifugal filtration apparatus is that the filtration rate starts out high because of its relatively large filter surface area. However, as the concentrate volume approaches its dead stop volume, the active filter area approaches zero. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume. Accordingly, this device is limited because it will filter to different dead stop volumes when used in centrifuge rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, and the filtration rate will slow dramatically as the concentrate volume approaches the dead stop volume.\nThe types of filtration devices disclosed in U.S. Pat. No. 5,112,484 to Zuk are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. In this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. Therefore, the dead stop volume will be different when the device is used in a swinging bucket rotor centrifuge than it will be when used in a fixed angle rotor, and will also be different when the device is used in fixed angle rotor centrifuges of different rotor angle. This means that the results obtained from this type of device when used in one type of centrifuge rotor cannot be compared to the results obtained from this type of device when used in another type of centrifuge rotor. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it must contain a venting means that vents to the atmosphere in the centrifuge). Another problem with this type of centrifugal filtration apparatus is that as the concentrate volume approaches its final retained volume, the active filter area approaches zero. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume. Although the retained concentrate can be removed from this device with a pipette, the retained concentrate is best removed from this type of device using the concentrate cup with a second spin. This type of device is not well suited to filter small volumes. Accordingly, this device is limited because it will filter to different dead stop volumes when used in centrifuge rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, the filtration rate will slow dramatically as the concentrate volume approaches the dead stop volume, it is not designed to filter small volumes, and it is not easy to remove the dead stop volume with a pipette.\nThe types of filtration devices disclosed in U.S. Pat. No. 5,490,927 to Herczeg are capable of preventing filtration to dryness and contain a dead stop feature which causes filtration to cease while there is concentrate remaining within the apparatus. In this type of filtration device, the amount of concentrate remaining after dead stop is dependent upon the angle of the centrifuge rotor. Therefore, the dead stop volume will be different when the device is used in a swinging bucket rotor centrifuge than it will be when used in a fixed angle rotor, and will also be different when the device is used in fixed angle rotor centrifuges of different rotor angle. This means that the results obtained from this type of device when used in one type of centrifuge rotor cannot be compared to the results obtained from this type of device when used in another type of centrifuge rotor. Another problem with this type centrifugal filtration apparatus is that is an open system (i.e. it must contain a venting means that vents to the atmosphere in the centrifuge). This type of device is not well suited to filter small volumes. This type of device utilizes the entire filter surface area until the dead stop volume is reached when used in a swinging bucket centrifuge. However when this type of device is used in a fixed angle rotor centrifuge the active filter area will approach zero as the concentrate volume approaches its final retained volume. Therefore, the filtration rate will slow dramatically as the concentrate volume approaches its final retained volume when used in a fixed angle centrifuge. This type of device requires a second spin to remove the retained concentrate solution. Accordingly, this device is limited because it will filter to different dead stop volumes when used in centrifuge rotors with different angles, it is an open system and will vent potentially harmful gases to atmosphere during filtration, the filtration rate will slow dramatically as the concentrate volume approaches the dead stop volume when used in a fixed angle rotor centrifuge, it is not designed to filter small volumes, and it requires a second spin to remove the retained concentrate solution.\nFinally, the art discloses the use of a lower retentate pocket to address the issue of filtering to dryness with centrifugal filtration devices. U.S. Pat. No. 5,647,990 to Vassarotti reveals a device having a concentrate retention pocket located at the bottom of a pre-filtration chamber, below the level of the filter membrane. Fluid flows through the filter into filtrate outlet channels, exiting these channels into the filtrate reservoir. These exit channels do not collect filtrate.\nIt is therefore an object of the present invention to provide a filtration device that may be used in a swinging bucket centrifuge, as well as a fixed angle centrifuge, without filtering to dryness.\nIt is also an object of the present invention to provide a filtration device that will filter to approximately the same dead stop volume when used in either a swinging bucket rotor centrifuge, or when used in a fixed angle rotor centrifuge.\nIt is also an object of the present invention to provide a filtration device having a relatively high filtration membrane surface area thereby enabling filtration to occur at a higher rate.\nIt is also an object of the present invention to provide a filtration device that continues to use a relatively high filtration membrane surface area until the dead stop volume is attained, thereby enabling filtration to occur at a higher rate throughout the entire filtration process.\nIt is also an object of the present invention to provide a filtration device, which minimizes the clogging of the semi-permeable membrane thereby maximizing filter efficiency.\nIt is also an object of the present invention to provide a filtration device from which the retained concentrate solution can be easily removed with a pipette.\nIt is also an object of the present invention to provide a filtration device which can be used as a closed system (i.e. does not vent to atmosphere), or that can be used as an open system (i.e. vents to atmosphere).\nIt is a further object of the present invention to provide a filtration device, which can be manufactured as a disposable small volume device, or as a disposable large volume device.\nThe foregoing problems of the prior art are solved, and the objects of the present invention are achieved, by use of a filtration apparatus constructed in accordance with the principles of the present invention.\nIn accordance with the present invention, the filtration apparatus for filtering and concentrating a solution includes a concentrate tube, which is divided into two sections. The upper cylindrical section acts primarily as a reservoir for unfiltered solution and directs unfiltered solution from the upper section into the lower section of the concentrate tube. The upper section is shaped so that all of the solution in the upper section will flow into the lower section regardless of the centrifuge rotor angle. The lower section of the concentrate tube is mainly rectangular, but may be triangular shaped at the bottom. The volume of the lower section is made as small as possible to insure that the entire filter surface area is used to filter the maximum amount of solution. The lower section may contain a pipette tip channel to facilitate the removal of the retained concentrate with a pipette. The lower section is open on one side, or two sides.\nA micro-porous membrane or membranes may be sealed to the lower section in several ways. First, the membrane can cover the open side of the lower section of the concentrate tube, or the membranes can cover the open sides of the lower section of the concentrate tube. Alternately the micro-porous membrane may be sealed to the filter cover, or covers. A second alternative is to seal the micro-porous membrane or membranes between the filter cover and a filter sealing gasket or gaskets. A third alternative is to seal the micro-porous membrane between the filter cover and the concentrate tube using a compression rib that compresses a portion of the outer periphery of the micro-porous filter between the filter cover and the concentrate tube.\nA filter cover is sealed to the open side of the lower portion of the concentrate tube, or two filter covers are sealed to the two open sides of the lower portion of the concentrate tube. The filter cover or covers contain a filter support means such as filter support ribs. The filter cover or covers also contains a filtrate outlet port. The position of the filtrate outlet port or ports determines the retained concentrate volume. The concentrate tube may contain a port or ports to allow displaced filtrate air to flow into the upper part of the concentrate tube to replace the filtered solution that flows from the concentrate tube into the filtrate tube.\nThe means for collecting the filtrate may include a filtrate tube. The concentrate tube may be sealed into the filtrate tube. The filtrate tube cap may contain a means to nest the concentrate storage tube."}
{"text": "In online gaming, game hosting services and game developers have created a number of ways to track and personalize the online gaming experience. One drawback of existing systems is that many of the features have grown up independent of each other. Games send blobs of data about gamers back and forth to a central service, but the service has no way to understand and aggregate the data outside of the game context. Games can host their own Websites, but the data displayed there is not universally accessible to other games.\nIn a sense, then, the service and games offer two parallel communities that offer great—but separated—resources for gamers. First, in the game community, while playing a game, the gamer can see the community of others who play the specific game, the leaderboards for that game, and his personal achievements in that game. A game can tell a gamer, from the Service data, if a Friend is online, but it can't tell the gamer what, exactly that Friend is doing on the Service or when he will be available.\nSecond, in the service community, the service knows a gamer player's history, all of the games he's played, the amount of time he spends online, the size of his Friends list and all of the games that Friends have played or are playing, the Friends invites sent and received, the Messages sent and received, and all of the Feedback the gamer has given and received.\nSystems have tried to leverage these on-line communities to match various players to allow them to play multi-player games. Nevertheless, in general such systems, which typically emphasize skill or experience in a single game or small family of games, do not allow groups of players who are likely to enjoy shared interaction based on a variety of personal considerations to form a group outside of the game context and then collectively engage in a multiplayer game.\nFurther, because matchmaking interfaces have typically been part of each individual game's code, gamers have not had the ability to create a match with other players in one game and then roam with that group of players to another game."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical actuator for controlling position of a focusing lens that focuses a laser beam on a disk surface. The optical actuator drives the focusing lens in three directions including focusing, tracking and tilting directions.\n2. Description of the Background Art\nAn Optical actuator drives a focusing lens which forms a beam spot on an information disk, such as a DVD (digital versatile disk) or a CD (compact disk). By driving a focusing lens in a focusing direction, which is orthogonal to the disk surface, and a tracking direction, which is parallel to a radial direction of the disk. The optical actuator controls position of the focusing lens, so that the beam spot is formed on a pit line of the disk accurately. Furthermore, when the disk surface tilts or wobbles with respect to an optical axis of the focusing lens due to disk warpage, the optical actuator performs tilt control by pivotally rotating the focusing lens on an axis along the tangential direction or on an axis-along the radial direction of the disk.\nFIG. 10 is a diagram illustrating a conventional optical actuator described in Japanese Patent Laid-Open No. 2001-297460. A focusing lens 101 is fixed to a lens holder 102. Six blade springs 103a–103c, and 103d–103f (103e and 103f are not shown) are fixed to each side of the lens holder 102. Print coils 104a and 104b are fixed to another side of the lens holder 102. A suspension holder 106 and permanent magnets 107a–107d are fixed to a base 105. The print coil 104a is disposed between the permanent magnets 107a and 107b, and the print coil 104b is disposed between the permanent magnets 107c and 107d. The blade springs 103a–103c and 103d–103f are connected to the suspension wires 108a–108c and 108d–108f (108e and 108f are not shown), each of which are fixed to the suspension holder 106.\nWhen a current is supplied to focusing coils built in the print coils 104a and 104b so as to generate electro-magnetic forces in a vertical condition, the lens holder 102 is controlled in the focusing direction Fo. When a current is supplied to tracking coils built in the print coils 104a and 104b so as to generate electro-magnetic forces in the radial direction, the lens holder 102 is controlled in a tracking direction Tk. When a current is supplied to the focusing coils so as to generate electro-magnetic forces having opposite direction in each of the print coils 104a and 104b, rotational moment on an axis in the tracking direction Tk occurs to the lens holder 102. As a result, the lens holder 102 is driven in the tilting direction Ti.\nWhen the lens holder 102 moves in the tilting direction, each of the blade springs 103a and 103c is warped or stretched, and the tilt control becomes unstable. Furthermore, the structure using blade springs requires higher cost and many assembly steps.\nIt is therefore, a primary object of the invention to provide an optical actuator which is able to perform the tilt control stably and can be produced with low cost."}
{"text": "Hospitals, clinics, and other healthcare providers continually face a need for donated blood to help in treating patients. It is often difficult to notify and educate the public regarding donations. Donors may not become repeat donors due to a lack of information and access. Operators and administrators may introduce inefficiencies in their operation and management of blood collection instruments due to a lack of information and access."}
{"text": "This invention relates in general to bird feeders and animal feeders, and is more particularly directed to demand feeders for agricultural birds, e.g., fowl such as chickens, ducks, turkeys, peafowl, or many other exotic or native birds, or small animals like cats and rabbits.\nThe invention is more particularly concerned with a feeder that dispenses feed on demand from a raised hopper onto a feed tray or catch dish that is suspended beneath the hopper, where the dispensing occurs by action of the bird or birds feeding, i.e., by action of the birds pecking and grabbing at food on the feed tray.\nIt is a significant object of the invention to create a bird feeder that is simple, easy to use, refill, and clean, and which allows for choice of a wide variety of feed types, from small grain, such as millet seed or mash, to larger particles such as pellets, corn, or dog food.\nThe bird feeder provides on-demand gravity feed and engages the birds in pecking and grabbing activities which helps stimulate normal, healthy behavior. The on-demand feeder also reduces over-feeding and eliminates food waste, dispensing feed only when needed."}
{"text": "Seagoing vessels are required in various types of offshore operations, including scientific surveys and oil and gas drilling and production. Such vessels are typically configured as drillships, barges and jack-up rigs, as well as supply and service ships.\nBoth active and passive techniques have been proposed for damping the natural oscillatory displacement of vessels to achieve roll and pitch stabilization in seas having periodic wave motion. Such techniques have utilized water tanks on the vessels and various arrangements of blowers, pumps, valves actuators, roll and pitch sensors, and electronic control circuits for moving water in the tanks to counteract oscillatory roll and pitch motion. A common objective of heretofore known systems has been to make the natural frequency of oscillatory flow of water in the tanks approximately the same as the natural frequency of oscillations of the vessel, thereby to \"tune\" the tanks to the vessel. Once tuned, the damping action is achieved by causing the flow pattern of water in the tanks to be approximately 90.degree. out of phrase with the natural oscillations of the vessel. The forces produced by the water in the tanks then tend to counteract the roll and pitch forces on the vessel.\nPrior systems that attempt to achieve stabilization in the manner described above have the disadvantage that large counteracting, damping forces must be produced in order for the system to be effective. The equipment required to provide the large counteracting forces is complex and expensive. In active systems, high power blowers and pumps are typically required. Passive systems generally require high-capacity valves, special stabilizing tank configurations and control circuits for timing the flow of water in the tanks."}
{"text": "Recently, digital discs have been widely used in many fields. For example, a CD (abbreviation of Compact Disc) is used for digital audio records. A CD-ROM (ROM; abbreviation of Read Only Memory) is used for data memories. An LV (abbreviation of Laser Vision) disc is used for video records. In those digital discs, recorded data or information are reproduced at a speed specified to each kind of the digital discs. Thus, reproduction apparatus of digital discs are designed to reproduce the digital discs at the sole reproduction speed, i.e., a standard speed specified to each of the digital discs.\nReferring now to FIG. 1, an example is shown of the conventional reproduction system for digital discs, e.g., a CD dubbing system. The CD reproduction system is provided for carrying out a dubbing from CD players to analog tape recorders.\nIn FIG. 1, the CD reproduction system comprises a CD player 10 and an analog tape recorder, e.g., a cassette tape recorder 11. The CD player 10 includes a CD reproducing mechanism section 12, a clock generator 13 and a signal processor 14.\nThe CD reproducing mechanism section 12 comprises a disc drive motor 15, a disc turntable 16 and a pickup device 17. The clock generator 13 generates a clock pulse Pck with a prescribed frequency Fck. The clock pulse Pck is applied to the signal processor 14.\nThe signal processor 14 comprises a servo control section 14a and an information signal reproducing section 14b. The servo control section 14a is coupled to the disc drive motor 15 and the pickup device 17 for supplying servo control signals thereto, as described later in brief. The information signal reproducing section 14b is coupled to the pickup device 17 for receiving an RF signal (RF; abbreviation of Radio Frequency) Srf read from a CD 18 and sent to the CD reproducing mechanism section 14.\nThe disc drive motor 15 drives the disc turntable 16 under the servo control of the servo control section 14a. The pickup device 17 shifts in the radial direction of the CD 18 under the servo control of the servo control section 14a. Thus, the pickup device 17 traces a record track on the CD 18 in a CLV (Constant Linear Velocity) speed. The pickup device 17 reads the CD 18 during the tracing of the record track of the CD 18 so that an RF signal Srf corresponding to the digital data recorded on the CD 18 is output from the pickup device 17.\nThe pickup device 17 outputs the RF signal Srf corresponding to the digital data, as described above. The RF signal Srf is applied to the information signal reproducing section 14b of the signal processor 14. A data slice circuit (not shown) in the information signal reproducing section 14b slices the RF signal Srf at a prescribed threshold level so that the RF signal Srf is shaped to the original digital data.\nThe servo control section 14a of the signal processor 14 generates two servo control data SCm and SCp and a bit synchronizing signal BS from the reproduced digital data. A motor servo control data SCm is applied to the disc drive motor 15, as described before. A pickup servo control data SCp is applied to the pickup device 17, as described before. The servo controls for both the disc drive motor 15 and the pickup device 17 are carried out in using the digital data obtained from the output of the pickup device 17, i.e., the RF signal Srf, and also the clock pulse Pck applied from the clock generator 13.\nThe information signal reproducing section 14b of the signal processor 14 demodulates an analog data AD from the reproduced digital data under the control of the clock pulse Pck. The analog data AD is applied to the cassette tape recorder 11. The cassette tape recorder 11 then records the analog data AD.\nThe conventional reproduction apparatus of digital discs, however, have a drawback that a reproduction of the digital discs takes a relatively long time the same as a standard record time provided for each of the digital discs."}
{"text": "1. Field of the Invention\nThe present invention is related to wireless tracking systems and methods. More specifically, the present invention relates to a system and method for mitigating multipath errors associated with the wireless tracking of objects by utilizing tags that transmit signals at various power levels.\n2. Description of the Related Art\nThe ability to quickly determine the location of objects located within a facility is becoming a necessity of life. To the uninformed observer, the placement of transponders, also known as tags, on numerous non-stationary objects whether in an office or home would appear to be an unnecessary use of resources. However, the uninformed observer fails to appreciate the complexity of modern life and the desire for efficiency, whether at the office or home.\nFor example, in a typical hospital there are numerous shifts of employees utilizing the same equipment. When a new shift arrives, the ability to quickly locate medical equipment not only results in a more efficient use of resources, but also can result in averting a medical emergency. Thus, the tracking of medical equipment in a hospital is becoming a standard practice.\nThe tracking of objects in other facilities is rapidly becoming a means of achieving greater efficiency. A typical radio frequency identification system includes at least multiple tagged objects each of which transmits a signal, multiple receivers for receiving the transmissions from the tagged objects, and a processing means for analyzing the transmissions to determine the locations of the tagged objects within a predetermined environment. The tags can also receive signals.\nSeveral prior art references discloses various tracking systems.\nMcKee et al., U.S. Pat. No. 6,915,135 discloses a system for determining presence, identity and duration of presence in a given area (a table in a restaurant) of an object (tag attached to a waiter).\nLester, U.S. Pat. No. 3,805,265 discloses a location system that uses line-of-sight radiant wave energy for signal transmission.\nSchwengler U.S. Pat. No. 7,050,819, is directed at the problem of adequate power for a mobile telephone for a two-way communication function or a regeneration function as a node of a mesh network.\nChrist, U.S. Pat. No. 5,977,913, discloses a radiofrequency system that is utilized within a prison and allows for an individual to be located after an alarm is triggered by the individual.\nZodnik, U.S. Patent Publication Number 2004/0147232, discloses wall-mounted (RJ-11 or RJ-45) wireless transceivers configured to only track the location of a self-identified wireless communication device in order to communicate the location of the self-identified wireless communication device to an emergency service such as 911.\nOne exemplary method triangulates the strongest received signals to determine the location of a tagged object. This method is based on the assumption that the receivers with the strongest received signals are the ones located closest to the tagged object. However, such an assumption is sometimes erroneous due to common environmental obstacles. Multipath effects can result in a further located receiver having a stronger received signal from a tagged object than a more proximate receiver to the tagged object, which can result in a mistaken location determination. The prior art has disclosed various means for overcoming multipath effects.\nTekinay, U.S. Pat. No. 6,259,894 for a Method For Improved Line-Of-Sight Signal Detection Using RF Model Parameters, discloses a method for reducing time-shift due to multipathing for a RF signal in an RF environment.\nClose, U.S. Pat. No. 3,869,673 for a Method And Apparatus For Measuring Multipath Distortion, discloses a method for indicating multipath distortion in a received signal.\nLennen, U.S. Pat. No. 5,402,450 for a Signal Timing Synchronizer, discloses a method and apparatus for reducing the effects of multipath induced distortions on the accuracy of detecting the time of arrival of a received signal.\nFortune et al., U.S. Pat. No. 5,450,615 for a Prediction Of Indoor Electromagnetic Wave Propagation For Wireless Indoor Systems, discloses techniques for predicting RF propagation within a structure.\nOther prior art references have disclosed the use of varying energy levels. For instance, Nakanishi, U.S. Pat. No. 5,451,847 for a Variable Energy Radio Frequency Quadrupole Linac discloses changing radio frequency energy levels to emitted focused and accelerated beams.\nKaewell, Jr. et al., U.S. Pat. No. 7,082,286 for a Path Searcher Using Reconfigurable Correlator Sets discloses producing a path profile for a user based on sorted output energy levels.\nFernabdez-Cobaton et al., U.S. Pat. No. 6,697,417 for a System And Method Of Estimating Earliest Arrival Of CDMA Forward And Reverse Link Signals discloses a mobile station receiver that detects the arrival times and energy levels of received signals, and constructs a searcher histogram and a finger histogram associated with each pilot signal.\nThe prior art has yet to resolve mistaken location calculations based on multipath effects for objects tracked within an indoor facility."}
{"text": "In recent years, airbag devices have become standard equipment for most vehicles. The airbag device is a safety device activated upon an emergency event such as a vehicle collision and inflated and deployed by gas pressure to receive and protect an occupant. There are different kinds of airbag devices for various installation locations or uses. For example, a front airbag is provided in the center of the steering at the driver seat in order to protect the occupant of the front seat against an impact mainly in the front-back direction, and a passenger airbag is provided in the vicinity of a front passenger seat at locations including the instrument panel. In addition, a side curtain airbag that is inflated and deployed along the side windows is provided at the side door area near the roof in order to protect occupants of the front and back seats against a side collision and a subsequent roll over (turn over).\nIn an occupant protecting device disclosed in Japanese Patent Application Publication No. 2014-034356, a side airbag to be inflated and deployed immediately next to an occupant is provided on each side of a seat. According to the above mentioned reference, occupant's movement to the side opposite to the collision side (one side in the vehicle widthwise direction) that may be caused by reaction against movement to the collision side (the other side in the vehicle widthwise direction) in a collision can be restricted.\n[PTL 1]\nThe airbag device described above is adapted to mainly cope with collisions in various directions such as frontal and side collisions, and airbags for a driver's seat and passenger seats are designed individually to cope with a frontal collision. For a side collision, separate airbags are used depending on the body parts to be protected by the airbags, for example the head is protected by a side curtain airbag, and the chest and abdomen are protected by side airbags. These airbags are supported by reaction force surfaces on the vehicle body side such as side surfaces of the vehicle interior and therefore have high occupant restraining capabilities even when the airbags receive the weight of an occupant in restraining the occupant. However, in order to cope with collisions in more varied directions, additional airbags are used. Such additional airbags may not always secure appropriate reaction force surfaces on the vehicle body side depending on where these airbags are installed, additional spaces for these airbags are not always easily available, and the vehicle weight and device cost might increase as a result.\nIn view of the problem, it is an object of the present invention to provide an occupant restraining device capable of coping with a plurality of collision directions without additional airbags."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nIt is known to provide a surgical jig having a cutting guide, which assists a surgeon in making accurate and repeatable cuts during surgery. The accuracy of a cut made using a cutting guide is mainly determined by the successful positioning and fastening of the cutting guide in the operative area. Typically, a stylus is attached to the cutting guide, which references from a reference point or marker on the body, for example, a bone. The stylus is either free to translate and/or rotate relative to the cutting guide, so that it may be adjusted, or is rigidly fixed to the cutting guide. An attachment mechanism mounts the stylus to the cutting guide, and a locking mechanism locks the stylus in a desired position relative to the attachment mechanism and hence the cutting guide. Once in the operating position, the jig can be attached to, for example, a bone using pins and/or bone screws.\nThe use of a stylus to position a cutting guide is common in joint replacement surgery. For example, during a total or partial knee replacement operation, it is necessary to resect the proximal tibia to a specific level. Existing tibial styli reference from the tibial plateau. This area of the joint is subject to considerable wear and, in a knee requiring joint replacement is likely to be damaged. Therefore, the tibial plateau does not provide a reliable reference position. It is desirable to provide a stylus assembly that enables referencing from an area of undamaged bone, which would provide a more reliable reference position.\nOnce a stylus has been used to position a cutting guide, the stylus is usually removed, in order to increase the working space available to the surgeon. A problem of existing arrangements is that the removal of the stylus can cause unintentional and unnoticed movement of the cutting guide. This leads to inaccurate positioning of the cutting guide, with the result that the surgeons' incisions are misdirected. It is therefore also desirable to provide a stylus assembly which can be removed with the minimum of handling."}
{"text": "The invention relates to an integrated circuit comprising at least one level of metallization, the level of metallization being provided with tracks and comprising metal portions of at least two different thicknesses.\nA level of metallization usually comprises metal tracks of small thickness, for example 1 micron, and small width, for example 1 micron. These metal tracks are used to transfer logic data. To manufacture inductors with an inductance of several nanohenrys, a separate level of metallization may be provided which is employed only for this purpose. In order to reduce the resistance of these inductors, particularly the direct-current resistance, use is made of inductors having a large width of, for example, 20 .mu.m, and a large thickness of, for example, 2.5 .mu.m.\nHowever, it is rather unsatisfactory to use a level of metallization only for inductors as a large area of this level of metallization remains unused. For this reason it would be desirable to take measures enabling said level of metallization to also comprise tracks for the transfer of logic data, said tracks having a small width. However, it is very difficult to manufacture tracks which are narrow and thick."}
{"text": "(1) Field of the Invention\nThe present invention relates to a device which releases a development affecting gas for affecting the development of plant materials and to a method for use of the device. In particular, the present invention relates to a method for initiating ripening or retarding ripening of a fruit or vegetable, particularly with ethylene.\n(2) Description of Related Art\nEthylene and other growth affecting gases have been used in storage rooms to modify the natural growth rate of plant materials including whole plants and plant parts such as fruits and vegetables. Such storage room systems are described in U.S. Pat. Nos. 3,620,765 to McDonnell et al; 3,924,010 to Erb; Re 28,995 to Burg. and 4,764,389 to LaBarge. These systems can only be used by wholesalers and others dealing in bulk fruits and vegetables.\nThe prior art has described various forms of packaging to inhibit or accelerate growth of fruits or vegetables. U.S. Pat. Nos. 3,069,274 to Concannon; 3,798,333 to Cummin; 4,528,228 to Clevenger and 4,622,229 to Toshitsugu are illustrative of such packaging. The common concept is the provision of a gas to facilitate ripening or the removal of a gas (ethylene) by the package to retard ripening of the fruit or vegetable. Such packaging can be expensive and the effect on development is subject to variation depending upon the handling and storage of the fruit or vegetable.\nThe prior art has described the treatment of plants to prevent fungal or other microorganism infection of cut or exposed portions. Illustrative is U.S. Pat. No. 2,923,632 to Castberg. Such treatments can only be used where the portion of the plant being treated is not eaten. Further, unless the fungicide or other chemical is safe for food purposes, it is difficult to obtain approval to market such devices. In any event, such devices are not used to affect the development (ripening or retarded ripening) of the plant material.\nThe prior art has described numerous chemicals for producing ripening in the manner of ethylene. They are particularly set forth by Kays, S. and Beaudry, R. in Acta Horticulturae 201 77-115 (1987). Burg, S. et al., Plant Physiol. 42 144-152 (1987) describes various unsaturated organic compounds as substitutes for ethylene to provide ripening of fruits and vegetables. Alternatively, carbon dioxide and other carbon-oxygen gases have been used to retard ripening of fruits and vegetables.\nThe problem has been to provide a convenient, on demand method and means for initiating the ripening and/or retardation of ripening of the fruit or vegetable."}
{"text": "I. Field of the Invention\nThe invention relates generally to the utilization of dimethyl ether and the production of hydrogen and carbon oxides. More specifically, the invention relates to a process for catalytically reacting dimethyl ether in the presence of steam.\nII. Description of the Prior Art\nThe conversion of low molecular weight alkanes, such as methane, to synthetic fuels or chemicals has received increasing attention as low molecular weight alkanes are generally available from secure and reliable sources. For example, natural gas wells and oil wells currently produce vast quantifies of methane. In addition, low molecular weight alkanes are generally present in coal deposits and may be formed during mining operations, in petroleum processes, and in the gasification or liquefaction of coal, tar sands, oil shale, and biomass.\nMany of these alkane sources are located in relatively remote areas, far from potential users. Accessibility is a major obstacle to effective and extensive use of remotely situated methane, ethane and natural gas. Costs associated with liquefying natural gas by compression or, alternatively, constructing and maintaining pipelines to transport natural gas to users are often prohibitive. Consequently, methods for converting low molecular weight alkanes to more easily transportable liquid fuels and chemical feedstocks are desired and a number of such methods have been reported.\nThe reported methods can be conveniently categorized as direct oxidation routes or as indirect syngas routes. The direct oxidative routes convert lower alkanes to products such as methanol, gasoline, and relatively higher molecular weight alkanes. In contrast, the indirect syngas routes involve the production of synthesis gas as an intermediate product.\nRoutes are known for converting methane to dimethyl ether. For example, methane is steam reformed to produce synthesis gas. Thereafter, dimethyl ether and methanol can be manufactured simultaneously from the synthesis gas, as described in U.S. Pat. No. 4,341,069 issued to Bell et al. The '069 Patent recommends a dimethyl ether synthesis catalyst having copper, zinc, and chromium co-precipitated on a gamma-alumina base. Interestingly, the '069 Patent states that dimethyl ether product can be charged as power generator fuel to a combustor of a gas turbine-prime mover arrangement, either alone or simultaneously with synthesis gas.\nAlternatively, methane is converted to methanol and dimethyl ether is subsequently manufactured from methanol by passing a mixed vapor containing methanol and water over an alumina catalyst, as described in an article by Hutchings in New Scientist (Jul. 3, 1986)35.\nHaving a relatively low vapor pressure, dimethyl ether is readily transportable. Moreover, dimethyl ether can be economically produced in relatively small quantities, as compared to materials such as compressed natural gas which require economies of scale associated with large cryogenic plants to be produced competitively. On the other hand, synthesis gas reportedly produces very little atmospheric pollution when combusted with air as fuel. Therefore, a practical process for converting dimethyl ether to synthesis gas on a commercial scale would be attractive to, among others, natural gas producers situated far from fuel consumers.\nKnown processes for producing synthesis gas typically react hydrocarbons with steam at elevated temperature over a catalyst. Generally, more complex hydrocarbons are converted to methane which is then steam reformed to produce hydrogen or synthesis gas.\nUnited Kingdom Patent Application GB 2213496 A listing Lywood as inventor describes the production of hydrogen-containing gas streams by an endothermic catalyzed reforming between methane and steam. The '496 Application proposes the following equations for the steam reforming of methane: EQU CH.sub.4 +H.sub.2 O.fwdarw.CO+3H.sub.2 1. EQU CH.sub.4 +2H.sub.2 O.fwdarw.CO.sub.2 +4H.sub.2 2. EQU CH.sub.4 +CO.sub.2 .fwdarw.2CO+2H.sub.2 3.\nU.S. Pat. No. 4,592,903 issued to Osman et al., states that carbon monoxide can be endothermically converted to carbon dioxide and hydrogen through a reaction termed a water-gas shift, represented by the equation: EQU CO+H.sub.2 O.fwdarw.CO.sub.2 +H.sub.2 4.\nReportedly, the \"shift\" reaction, can be accomplished in two shift conversion vessels operating at different temperatures to maximize yield. The '903 patent states that a temperature of from about 600 to 900 degrees F. and a pressure of about 300 to 1,000 psig is effective in a high-temperature shift converter containing a supported, chromium-promoted iron catalyst. The '903 Patent further states that a low-temperature shift conversion takes place over a catalyst comprising a mixture of zinc and copper oxides at a temperature of from about 400 to 500 degrees F. and a pressure of from about 300 to about 1,000 psig.\nIt is important to distinguish between the steam reforming of hydrocarbons, as described above, and the partial oxidation of hydrocarbons. The partial oxidation of methane produces two moles of diatomic hydrogen for each mole of methane reacted. In contrast, the steam reforming of methane produces three moles of diatomic hydrogen per mole of reacted methane.\nThe partial oxidation of methane is described, for example, in U.S. Pat. No. 4,618,451 issued to Gent. The '451 Patent states that methane is reacted with oxygen from an air separation plant, the proportion of oxygen being less than sufficient for complete combustion. A hot gas containing hydrogen and carbon monoxide is said to be produced. The '451 patent also states that steam or nitrogen can be present during the combustion to act as a temperature modifier and to avoid soot formation. Additional hydrocarbon is reportedly injected into the hot gas, and the resulting gas mixture is reacted over a steam reforming catalyst.\nA particular class of partial oxidation processes for converting methane or natural gas to synthesis gas are known as autothermal processes. By convention, the autothermal process includes an exothermic oxidation step and an endothermic steam reforming step which are in approximate heat balance. For example, U.S. Pat. No. 5,112,257 issued to Kobylinski and assigned to the assignee of the present invention, describes an autothermal process for converting natural gas to synthesis gas which includes the steps of mixing natural gas with air, subjecting a resulting mixture to simultaneous partial oxidation and steam reforming reactions, and subsequently reacting unconverted alkanes with water in the presence of a catalyst having steam reforming activity.\nProcesses which produce hydrogen or hydrogen-containing mixtures by reacting a single-carbon saturated alcohol, methanol, with steam are collectively termed methanol steam reforming processes. U.S. Pat. No. 4,091,086 issued to Hindin et al. describes a process for producing hydrogen by reacting steam with methanol in the presence of a catalytic composition at elevated temperatures. The '086 Patent reports states that methanol can be converted to hydrogen in a single-stage reaction over a catalytic composition comprising zinc oxide, copper oxide, thorium oxide, and aluminum oxide. Moreover, the '086 Patent states, without citing authority or presenting evidence in support, that the composition catalyzes a purported methanol decomposition. The purported decomposition is described as producing significant amounts of carbon monoxide which are immediately consumed in a water gas shift reaction.\nU.S. Pat. No. 4,743,576 issued to Schneider et al. describes a catalyst for the production of synthesis gas or hydrogen from aqueous methanol by dissociation or steam reforming. The catalyst reportedly contains a noble metal component on an oxide carrier which comprises an oxide of cerium or titanium and, also, an oxide of zirconium or lanthanum.\nU.S. Pat. No. 4,865,624 issued to Okada describes a process for reacting methanol with steam including a decomposition reaction zone regulated at a temperature between 250 and 300 degrees C. and a conversion reaction zone regulated between 150 and 200 degrees. The '624 patent postulates an alleged methanol decomposition for producing hydrogen and carbon monoxide directly from methanol. The conversion reaction zone described in the '624 Patent is apparently intended to promote the well-known water gas shift reaction.\nAn integrated turbo-electric power generation system which incorporates methanol reforming as a source of fuel and as a means of heat recovery is described in sales literature circulated by the New Energy and Industrial Technology Development Organization under the authority of the Ministry of International Trade and Industry of Japan cerca 1985. The sales literature states that methanol and steam are passed through catalysts at temperatures in the range of 250 to 350 degrees C. to produce hydrogen and carbon dioxide in an endothermic reaction. The hydrogen-containing gas is reportedly combusted with air to drive a turbine. The sales literature indicates that the reforming reactor charge and the combustion air stream can be heat exchanged with the turbine exhaust to promote energy efficiency.\nDespite some earlier speculation regarding the existence of a direct methanol decomposition mechanism, practitioners generally agree that methanol steam reforming proceeds by a mechanism which does not involve the direct decomposition of methanol to hydrogen and carbon monoxide. Rather, it is accepted that the steam reforming of methanol creates methyl format and formic add as intermediaries. For example, an article by Jiang et al., Applied Catalysis A: General, 97 (1993) 145-158 Elsevier Science Publishers B. V., Amsterdam, cites studies and presents experimental data indicating that steam reforming of methanol proceeds via dehydrogenation to methyl formate, hydrolysis of methyl formate to formic acid, and decomposition of formic acid to carbon dioxide and hydrogen. According to the Jiang et al. article, no carbon monoxide production was detected while passing methanol over a copper, zinc oxide and alumina catalyst at temperatures below 250 degrees C. The Jiang et al. article reports that significant amounts of carbon monoxide were formed only at temperatures over 300 degrees C. Moreover, the Jiang et al. article states that methanol steam reforming proceeds in accord with the following equations: EQU 2CH.sub.3 OH.fwdarw.CH.sub.3 OCHO+2H.sub.2 5. EQU CH.sub.3 OCHO+H.sub.2 O.fwdarw.CH.sub.3 OH+HCOOH 6. EQU HCOOH.fwdarw.CO.sub.2 +H.sub.2 7.\nUK Patent Application GB 2085314 A listing Twigg as inventor describes a catalytic process for reacting a hydrocarbon with steam in net endothermic conditions to produce a gas containing carbon oxides and hydrogen. The process is reportedly carried out using a catalyst comprising the product of thermally decomposing and reducing intimately associated compounds of nickel and/or cobalt and at least one difficultly reducible metal. Reportedly, the catalyst also comprises a water-insoluble compound of an alkali metal oxide with an acidic or amphoteric oxide or mixed oxide.\nThe '314 Application states that the alkali metal, usually sodium or potassium, is chosen on the basis of the vapor pressure of its hydroxide form, so as to be available as an alkaline metal hydroxide to catalyze a reaction between carbon deposited on the catalyst and steam. The '314 Application speculates that the starting hydrocarbon can be any of those proposed for use with a catalytic steam/hydrocarbon reaction including methane, natural gas, liquified petroleum gas, naphtha, methanol, dimethyl ether, and isobutyraldehyde. However, as explained below the presence of potassium hydroxide actually hinders the reaction of dimethyl ether and steam.\nIn order to better utilize remotely situated sources of natural gas, to transport the energy inherent in natural gas in a safer and more economic manner, and to provide a fuel which creates very little atmospheric pollution when combusted in air, a commercially practical method for transforming dimethyl ether and steam to synthesis gas is desired. Preferably, the improved method is suitable for integration into modem power generation schemes."}
{"text": "A conventional fixing structure of a door weather strip for a car is described with reference to FIGS. 1 to 3. In some cases, insert clips 30 have been used for fixing a molding part 11 of a weather strip 10 to a door panel 1 of a car, which forms a corner part of the weather strip 10.\nThe insert clip 30 has a clip part 31 protruded from the center thereof so as to fix the weather strip 10 to the door panel 1, and a catch part 32 provided at the upper portion of the clip part 31 for preventing the insert clip 30 from coming off from the weather strip 10 or from rotating.\nHowever, in the conventional fixing structure of a door weather strip, since the molding part 11 of the weather strip 10 is fixed to the door panel 1 with the insert clip 30, the fixed structure of the insert clip 30 and the weather strip 10 is incomplete that the insert clip 30 and the weather strip 10 are not firmly integrated with each other.\nAccordingly, there arises a problem that a gap 33 is produced between the insert clip 30 and the weather strip 10 thereby lowering a sealing property between the weather strip 10 and the door panel 1. Further, there arises another problem that the insert clip 30 moves from its proper position in the weather strip 10 thereby lowering an external appearance."}
{"text": "1. Field of the Invention\nThis invention relates to plastic containers for rubbish, and the like, and is directed more particularly to a wheeled container having improved wheel mounting means.\n2. Description of the Prior Art\nWheeled containers for rubbish, refuse, and the like, are generally well-known. See, for example, U.S. Pat. No. 1,014,475, issued January 1912, to Chester L. Holloway; U.S. Pat. No. 3,366,397, issued Jan. 30, 1968, to Charles F. Zeilstra et al; U.S. Pat. No. 4,351,539, issued Sep. 28, 1982, to Michael S. Rodolakis; U.S. Design Pat. No. 218,359, issued Aug. 11, 1970, to William J. Marsh; and U.S. Design Pat. No. 231,184, issued Apr. 9, 1974, to Thomas E. Brown et al.\nSuch containers typically are molded of plastics material. Wheel assemblies for such containers often include two wheels mounted on a single rigid axle, usually of metal. The container may be molded with axle holes therein or axle holes may be cut in the container after molding. The axle usually is inserted through the two holes. One wheel may be fixed to the axle before attachment of the axle to the container but the remaining wheel must be fixed to the axle after the axle is in place. Alternatively, there may be molded in a bottom surface of the container a groove for receiving the axle. In such case, an axle with both wheels fixed thereto may be placed in the groove. The groove is then closed, at least in part, by a bracket, or the like, fixed to the container by fasteners.\nThere is a need for such a container having a wheel assembly which may be pre-assembled and mounted, as is, on the container quickly and easily without the need to bore holes in the container, attach wheels to axles (other than during pre-assembly of the wheel assemblies), or attach axle-retaining brackets to the container."}
{"text": "Semiconductor power devices have conventionally become the focus of attention, which are mainly used for systems in various power electronics fields such as motor control systems and power conversion systems.\nAs semiconductor power devices of this type, SiC semiconductor devices having a trench-gate structure have been proposed, for example.\nFor example, Patent Literature 1 discloses a field effect transistor including an n+-type SiC substrate, an n−-type epitaxial layer (drift region) formed on the SiC substrate, a p-type body region formed at a surface side of the epitaxial layer, an n+-type source region formed at a surface side within the body region, a grid-shaped gate trench formed in a manner penetrating through the source region and the body region to reach the drift region, a gate insulating film formed on the inner surface of the gate trench, a gate electrode embedded in the gate trench, a source trench formed in a manner penetrating through the source region and the body region to reach the drift region at a position surrounded by the grid-shaped gate trench, and a source electrode formed in a manner entering the source trench"}
{"text": "One of the many ailments of mankind is the formation of calculi, that is stones or concrements, within the body. More common places for the development of such concrements are in the renal system and the gallbladder. Some of these stones pass through the renal system without causing any adverse affects noticed by the person having the stones. Other stones lodge within the kidney, the bladders or passageways between kidney, the bladder and the exterior of the body causing pain or impairing the operation of organs; then the removal of the calculi becomes imperative.\nOriginally the only way of removing these stones was through invasive surgery. With the advent of improved imaging techniques it became possible to utilize percutaneous techniques for the removal of many of the calculi. The percutaneous techniques are also invasive and require the insertion of means such as a nephrotomy catheter to accomplish the stone removal. The catheters are equipped with devices to grasp the calculi or ultrasonic generators for fragmenting the calculi by vibrating while contigious thereto.\nMore recently, extra corporeal shock wave lithotripsy has been applied to fragment and thereby eliminate renal calculi.\nA lithotripsy system now being used has a pair of substantially orthogonally positioned X-ray generators designed to exactly locate the stone to be eliminated. The patient is moved three-dimensionally in a bath until cross-hairs on each of two images displayed by the system are juxtaposed to the stone. The cross-hairs are each at a line projection of the focal point of the shock wave generator. The patient is manipulated until the cross-hairs in each of the images are superimposed on the image of the stone; at that time the shock wave generator is triggered and the shock waves focused to shatter the stone down to sizes where the fragments will be washed out of the body through normal body functions.\nThe imaging systems used presently for the calculi locating function leave room for improvement. Among other things it is desirable to reduce the exposure of the patient to X-rays. The patient presently is exposed to X-rays, originally to ascertain that stones are indeed causing problems and that they can be eliminated by lithotripsy.\nX-rays are used to align the stone with the focal point of the shock wave generator and X-rays are used subsequent to the application of the shock waves to assure that the stones have been shattered into fragments which can be disposed of by normal functions of the body. Thus, the procedure requires a relatively large accumulated dose of X-rays. It would be advantageous to lower the amount of X-ray radiation required in the lithotripsy process. To this end ultrasound has been tried in the past. Additionally ultrasound radiation is more capable of imaging gallbladder stones than X-ray radiation. However, the prior art attempts at imaging with ultrasound for locating calculi to be fragmented by shock waves have encountered difficulties.\nIn the prior attempts at using ultrasound as the medium for locating the calculi, different acoustic windows were used by the ultrasound transducers and by the shock wave generator. The prior systems also failed to make any effective correction for the distortions of the ultrasound signals and the shock waves due to refraction caused by the differences of the transmission velocities of the ultrasound signals and the shock waves in the different media.\nIn addition in the past no provisions were made for enabling a single ultrasound transducer to image a plurality of planes for assuring accurate location of the calculi.\nSimilarly, no provisions were made for varying the field of view of the individual transducers to provide for the different views required when searching for calculi and when centering the discovered calculi at the cross-hairs.\nAccordingly it is an object of the present invention to overcome the shortcomings of the prior art ultrasound imaging systems used in calculi location and to thereby improve on the presently available calculi locating imaging systems used in cooperation with extra corporeal shock wave lithotripsy."}
{"text": "In many industries it is necessary or desirable to sequentially expose individual articles to specialized controlled environments. For example, in the food packaging industry several methods and devices exist for sequentially exposing containers of food products to a vacuum or to an inert atmosphere, and for sealing the container to retain the applied vacuum or atmosphere and thereby preserve freshness of the food products.\nOne technique for supplying a controlled environment utilizes individual \"bell jar\" enclosures which surround the container. Such a device is disclosed, for example, by U.S. Pat. No. 2,292,887 issued to McBean. Although such devices may be adapted for continuous operation, the apparatus necessary for moving the container and/or the \"bell jar\" relative to one another adds considerable complexity and cost to the device.\nOther devices are known which cooperate directly with the opening of the container to apply, for example, a vacuum substantially to the interior of the container only. U.S. Pat. No. 2,457,690 issued to Kronquest discloses an example of such a vacuum-head apparatus. Such devices, however, are complex, and are not readily adaptable to use with containers having differing configurations and sizes.\nIt is further known to utilize a multi-chambered rotary drum, in which individual chambers may accept a container to be subjected to the controlled environment or vacuum. As the drum rotates, the chambers transporting the containers move past and engage and outer enclosure which seals the opening of the chamber, and a vacuum is applied. The container is then sealed, normal atmosphere is returned to the chamber, and the drum rotates past the enclosure so that the sealed container may be removed. U.S. Pat. Nos. 1,751,643 issued to Malmquist, 2,521,746 issued to Preis, and 1,774,529 issued to Sharp illustrate various configurations of such known devices.\nUnfortunately, previous rotary drum devices have suffered from several shortcomings which limit their usefulness. For example, tolerances between the surface of the rotary drum and the outer enclosure must be very close in order to prevent excess leakage between the outer atmosphere and the enclosed inner chamber and controlled environment. The interface between the moving drum and the stationary enclosure is subject to wear and introduces considerable friction which must be overcome to rotate the drum. These problems are increased when a vacuum is applied to the chambers of the drum, for the pressure of the atmosphere on the outside of the stationary enclosure increases the frictional contact with the rotary drum. Such pressure may also distort the enclosure and interfere with a good vacuum seal.\nVarious configurations of rotary drums and enclosures have been utilized in attempts to minimize these difficulties, including specialized gaskets, sliding seals, and bearings described in the prior art. However, such attempts have been largely unsuccessful, and have added to expense and complexity which reduces reliability.\nFurther, prior systems have not been readily adaptable to accommodate numerous configurations and sizes of containers. As a result, it has been difficult or impossible to utilize a single device in conjunction with differing sizes or shapes of containers without extended down-time for reconfiguration of the system.\nAccordingly, it is an object of the present invention to provide a rotary drum sealing device and method which provide a secure closure of desired individual chambers during rotation of the drum without significantly interfering with such rotation. A related object is to provide such a method and apparatus which do not introduce significant friction between the sealing enclosure and the rotating drum.\nA further object is to provide a rotating drum vacuum sealing apparatus and method which eliminate the need for sliding seals.\nAn important object of the present invention is to provide a rotary drum vacuum sealing apparatus and method for use in a continuous sealing operation, which is mechanically simple, having few components, and which is therefore economical and highly reliable.\nAnother object is to provide such a sealing apparatus which is readily adaptable for use with varying sizes and configurations of jars and other containers. A related object is to provide such a device including means for closing the container which is adaptable to differing heights of containers. A further related object is to provide such a system adapted for receiving containers of differing diameters.\nThese and other objects shall be apparent in light of the present specification."}
{"text": "The present invention relates to apparatus, including a cartridge adapted to be removably mounted to the apparatus, for severing a tubular workpiece.\nTube cutting apparatus employed to sever sections of a continuous length of tubing or similar elongated stock fed from a source thereof such as a mill is known. Tube cutting apparatus in which a cut-off die set, driven by a ram assembly and clamped to the continuous length of tubing prior to and during the severing operation is exemplified by U.S. Pat. 4,653,368 issuing Mar. 31, 1987 to Riera et al., the disclosure of which being incorporated by reference herein. Riera discloses a tube cut-off die set including mechanism for providing a horizontal notching cut in a tubing wall followed by a vertical severing cut as the tube is fed through the set. The tub cut-off die set includes a movable upper shoe, a stationary lower shoe, guide posts and associated bearings affixed to the shoes and a ram providing the relative reciprocation between shoes. A tube holding arrangement therein comprises two axially spaced pairs of complementary die jaws being mounted on the lower die show for releasably clamping the tube. Downward reciprocation of a cam driver extending from the upper die shoe engages a pair of cam rollers associated with the lower die shoe to drive the die jaws about the tube and the blades between the jaw members.\nThe die set is used to cut sections of tubing as the tube is continuously manufactured. While suitable for its intended operation, the above-referenced die set does not allow for rapid adjustment of the die jaws so as to change their grip about the tube, and does not permit rapid replacement of the die jaws to accommodate different tube sizes and shapes. Since each tube diameter and cross section requires different jaw members it would be desirable to have jaw members that can be replaced and preset at a location spaced from the die set so that there would be no downtime of the tube-cutting apparatus. Further, a more compact jaw mounting arrangement would be desirable to reduce the overall size of the die set.\nIt is an object of this invention to provide a tube cut-off machine that maximizes the ease with which die jaws are replaced, to accommodate different tube diameters, without stopping the machine operation, and provide the user with an adjustment arrangement whereby the clamping force by the die jaws on the tube can be easily adjusted. Advantageously the use of a self-contained cartridge that includes the jaw members, cam followers and adjustment arrangement would allow the user to keep the machine cutting one particular diameter size while other cartridges are retrofitted with the appropriate jaw members, thereby reducing downtime of expensive machinery when it is time to change the jaws.\nPart of the setup time of jaw members associated with the Riera machine was removing the cam driver prior to removal of jaw holders and then remounting the jaw holders and cam driver. This arrangement has an element of machine downtime. It would be desirable to provide rapid connecting and disconnecting arrangement for the cam driver. Desirably, such an arrangement would associate the desired jaw member spacing with the cam driver.\nA further object of this invention is to provide a replaceable kit for use in a tube cutting machine which is simple, has jaw holders and jaw members preset, and is insertable as a unit. Advantageously, because the costs associated with the downtime of a machine is high and the costs of a cam driver relatively small, teardown and setup costs can be more than justified by such a kit. In this regard the kit advantageously permits only minimal downtime between removing cartridges and associated jaw holders when different tube diameters are to be cut.\nOne problem sometimes experienced in connection with tube severing machine is that the severed tube end portion is locally deformed and has what are termed \"dimples.\" Tubes that are so formed are waste and thus costly. While the Riera construction including a horizontal notching blade cantilevered at the end of a cross slide in connection with the vertical severing blade has worked well in most situations, tube dimpling with some cross sections and tube thicknesses results. A tube severing arrangement which does not produce deformed cross sections would be desirable.\nAnother object of this invention is provision of a tube severing arrangement that permits severing of thicker and larger diametered tubes without breaking the blades or increasing the size of the machine into which they are mounted. Further, configuring the notching blade so as to have an offset portion surprisingly allows such tube cross sections to be cut.\nThe foregoing objects and advantages are accomplished by a tube cut-off machine for severing elongated material moving continuously longitudinally of its length, comprising relatively reciprocable upper and lower die shoes, one shoe mounting a tube cut-off blade and a cam driver, and tube clamping means for releasably clamping the tube during the severing operation. The tube clamping means includes two laterally spaced jaw holders each carrying a jaw member and positioning the tube clamping surfaces of the jaw members in confronting relation and a cam follower operably connected to one jaw holder.\nIn accordance with this invention, the tube clamping mans is characterized by a compact, self-contained, cartridge that is removably mounted to the die set as a single by replaceable unit. The cartridge comprises a carriage adapted to be mounted on the mounting surface of the lower die shoe so as to be stationary relative thereto, a pair of laterally spaced cam followers each being mounted on a cam holder one and the other follower being a roller and a replaceable wear pad, the two jaw holders and cam follower holders, respectively, being laterally spaced and slidably disposed on the carriage in confronting relation, bias means for laterally biasing the jaw holders apart, and adjustment means including a frame having an adjustment screw and operably associated with the carriage for adjusting the relative lateral separation between the jaw members, the frame maintaining the holders in side-by-side relation and the screw adjusting wear in the wear pad.\nAdvantageously the holders are separately manufactured and thus the assembly is less expensive because the need for a \"beefed up\" carriage construction is minimized. Separate parts weigh less, thus reducing forces needed during clamping/unclamping of the tube to overcome friction and inertia. Advantageously, the wear plate absorbs shock loads which might otherwise shear the pivot pin supporting a cam roller for rotation, resulting in costly repair and machine downtime.\nA kit adapted to be mounted on the die set comprises the cam driver being operably disposed within the cartridge between the cam followers and their spacing preadjusted so that the jaw members are at their desired separation for clamping about the tube. The kit is mounted separately to the machine and requires minimal assembly time.\nIn accordance with one mounting arrangement, the cartridge is slid beneath a pair of gibs to secure the cartridge to the lower shoe, a dowel pin in a driver holder is connected to the upper shoe, and a pair of driver clamping plates on the driver holder are tightened behind the cam driver.\nThe cam driver has compliant load bars adapted to engage both cam followers to reduce uneven loads placed on the cam followers.\nFurther in accordance with this invention, the die set includes four bushings with their vertical axes disposed in a rectangular array with first, second and third adjacent pairs being secured together, respectively, by a first tie bar, a second tie bar and rail mounting plate, and the horizontal notching blade is vertically adjustable relative to a holding bar adapted to move between vertical planes including the first and third pair of bushings. Advantageously, a stiffening arrangement including three guide bushings including the second tie bar and rail mounting plate, the bashing axes forming a right triangular array, can stiffen the die set to obviate wobble of the notching blade."}
{"text": "Conventionally, a technique has been suggested for allowing an electromagnetic field to exist in an interspace area sandwiched between conductive sheets that face each other and for allowing the electromagnetic field to travel in a desired direction by changing the electromagnetic field by changing the voltage between the two conductive sheets or by changing the voltage between the conductive sheets by changing the electromagnetic field so as to transmit an electromagnetic wave. Non-patent document No. 1 shown in the following introduces an electromagnetic wave interface device that transmits/receives signals and receives electric power by electromagnetic wave transmission.\nAlso, a signal transmission system combined with a signal transmission device is known that transmits signals by changing an electromagnetic field in an interspace area sandwiched between a mesh-like conductor portion and a sheet-like conductor portion and in a transudation area on the outside of the mesh-like conductor portion.\nIn an electromagnetic wave transmission sheet used for the signal transmission system, the electromagnetic wave intensity in the transudation area is reduced according to the distance from the sheet. For example, the following non-patent document No. 1 discloses that arranging a resistor or an electromagnetic wave absorber at the end of the mesh-like conductor portion of the electromagnetic wave transmission sheet prevents the leakage or reflection of an electromagnetic wave.\n[Patent document No. 1] JP 2007-281678\n[Non-patent document No. 1] Hiroyuki Shinoda, et al., Method of Simultaneous Signal-Power Transmission Using Surface Microwave (Fundamental theories for ubiquitous and sensor networks, etc.), The Technical Report of The Proceeding of The Institute of Electronics, Information and communication Engineers, Vol. 107, No. 53(20070517) pp. 115-118"}
{"text": "In a conventional large-scale computer system that includes multiple chips, an error that occurs in the hardware that is contained in the computer system is detected and information on the detected error and information on the time at which the error occurred are stored in a register that is a storage device. A system managing apparatus, managing the computer system, accesses a register in which information on the error is stored and analyzes the error.\nA configuration of the above-described computer system will be described below with reference to FIG. 9. FIG. 9 is a diagram of a configuration example of the system according to the conventional technology. The system illustrated in FIG. 9 includes system boards on each of which central processing units (CPUs) are mounted, a crossbar board on which crossbar chips each including a timer are mounted, I/O boards on each of which a large scale integrated (LSI) circuit that controls the I/O device is mounted, and a system management board (MMB) that is a service processors (SVP) and on which firmware that manages the system is mounted. The system management board accesses the register of each chip via an inter-integrated circuit (I2C) bus.\nFor example, once a slave address and a channel of an I2C multiplexer that is an upper-level device with respect to a chip that is to be accessed are specified, the system management board is connected to the board on which the chip to be accessed is mounted. The system management board then specifies the slave address of the chip to be accessed and starts communications. The system management board then refers to the error information and the timer value (time information) that are stored in the register of the accessed chip and analyzes, for example, how the error occurred.\nHowever, as illustrated in FIG. 9, in a system in which broadcast communications cannot be performed, it is difficult to analyze an error that has occurred in the hardware. Specifically, in a system configuration in which broadcast communications cannot be performed, because the timers of the chips are not synchronized, it is not possible to analyze how an error occurred according to the timer value that is stored in the register and is referred to on the occurrence of the error.\nTo synchronize the timers (time information) in the above-described large-scale system, various technologies are disclosed in, for example, Japanese Laid-open Patent Publication No. 10-28119 and Japanese Laid-open Patent Publication No. 04-182806.\nThe above-described conventional technologies, however, have a problem in which it takes time to complete a process for synchronizing the timers (time information)."}
{"text": "1. Technical Field\nThe present invention relates to an electrostatic image developing toner, a method manufacturing for the electrostatic image developing toner, an electrostatic image developer, a toner cartridge, a process cartridge, an image-forming method, and an image-forming process.\n2. Related Art\nIn an electrostatic image developing toner using addition polymerizable resins having a random monomer chain and a polycondensation resin as a binder resin, fixation has been mainly accelerated by heating rather than pressure.\nIn recent years, in printing and copying techniques by electrophotographic methods, the industry has come to grips with the object of utilizing fixation by pressure (pressure fixing) to shift the course from conventional fixing method primarily by heat energy for the purpose of further energy saving."}
{"text": "Biodegradable softeners have attracted recent attention in the prior art. For instance, in German patent no. 197 43 687, in the name of Henkel KGaA, readily biodegradable detergents are described, which contain oligomeric esterquats obtained by quaternizing oligoesters of mono and di-carboxylic acids in combination with alkylene oxide adducts on fatty acid amines.\nThe international patent application WO-A-01/47489, in the name of Cognis Deutschland GmbH and Bigorra Llosas et al., discloses fiber brightening and softening agents comprising esterquats obtained by reacting alkanol amines with mixtures of fatty acids and dicarboxylic acids, optionally alkoxylating the resulting esters, and quaternizing the products; and auxiliary materials selected from (non)quaternized fatty acid amides, betaines, nonionic surfactants, polyols and/or their derivatives, alcohols and/or hydrotropes.\nIn the European patent application 1 136 471, in the name of KAO Corporation S.A., alkanol amine esters are described which are based on the esterification reaction of alkanolamines, carboxylic acids and fatty alcohols. The alkanolamines and fatty alcohols are optionally alkoxylated. In addition, the cationic surfactants and esterquats obtainable therefrom are disclosed.\nThe cationic surfactants and the esterquats disclosed in said European patent application have a high degree of biodegradability, but compared to the biodegradable esterquats of the prior art also exhibit a high degree of efficacy in softening and conditioning natural and synthetic fibers, such as hair, or fibers used in textiles and paper.\nIn a further aspect, said European patent application relates to aqueous fabric-softening compositions which contain the cationic surfactants or esterquats, optionally together with other active softening substances. Particularly, these fabric softening compositions contain, in an aqueous medium optionally containing constituents selected from those normally used in fabric softener composition: (a) cationic surfactants or esterquats obtainable from the alkanolamines described, (b) one or more cationic surfactants which are active as fabric softeners, and (c) one or more non-ionic fabric-conditioning surfactants, wherein the amount of (a)+(b)+(c) is 2-60 wt. % based on the total composition; the amount of (a), based on the total of (a)+(b)+(c), being 2-100 wt. %; the amount of (b), based on the total of (a)+(b)+(c), being 0-98 wt. %; and the amount of (c), based on the total of (a)+(b)+(c), being 0-40 wt. %."}
{"text": "Temperature-Compensated Xtal (crystal) Oscillator (hereinafter referred as TCXO) is a kind of crystal oscillator which can work in a wide temperature range and keep the output frequency of the crystal oscillator within a certain accuracy range (106˜10−7 orders of magnitude) through a certain compensation. It has a characteristic of low power consumption, working upon power-up, high stability and so on. Therefore it has been widely used in various communications, navigation, radar, satellite positioning system, mobile communication, program-controlled telephone switch and various electronic measuring instruments.\nThe temperature-compensated crystal oscillator in prior art is essentially a Voltage-Controlled Xtal (crystal) Oscillator (hereinafter referred as VCXO) with a temperature compensated network which produces a temperature-dependent compensation voltage. The key component in the uncompensated voltage-controlled crystal oscillator is a AT-cut quartz crystal, which temperature characteristic curve is approximately a cubic curve, and the cubic curve can be expressed as:f(T)=a3(T−T0)3+a1(T−T0)+a0  (1)\nWhere a3 is the coefficient of cubic term, a1 is the coefficient of linear term, a0 is the oscillating frequency at a reference temperature T0, T is the temperature of the location the AT-cut quartz crystal is at.\nThe frequency linear gain characteristics of the VCXO in prior art can be approximately expressed as:f(VC)=−G(VC−VC0)+f0  (2)\nwhere G is the gain of the VCXO, VC is the control voltage of the VCXO, VC0 is the input voltage of the voltage control terminal of the VCXO, f0 is the oscillating frequency when the input voltage is VC0.\nThen, the equation of the compensation voltage of the crystal temperature characteristic can be expressed as:VC(T)=A3(T−T0)3+A1(T−T0)+A0  (3)\nWhere A3=a3/G, A1=a1/G, A0 is a compensation voltage when the temperature T is T0.\nIn order to implement equation (3), it is necessary to generate a temperature compensation voltage applied to the VCXO for temperature compensation to counteract the frequency temperature characteristic, thereby obtaining a stable frequency output within a wide temperature range, and realizing the purpose of temperature compensation.\nFIG. 1 is a diagram of a TCXO based on digital circuit in prior art. The TCXO uses digital temperature compensation with an open-loop architecture. As shown in FIG. 1, it comprises a temperature sensor and conditioning circuit 101, a microprocessor 102, a compensation network 103, and a VCXO 104. The temperature T is acquired by the temperature sensor and conditioning circuit 101, and then input into the microprocessor 102 to find the compensation voltage value according to the temperature in the temperature-compensated voltmeter. Then the compensation network 103 converts the compensation voltage value into a compensation voltage, and inputs the compensation voltage to the voltage control terminal of VCXO 104 to make it output a signal with stable frequency. It can be seen that the temperature-dependent compensation voltage is directly applied to the voltage control terminal of the VCXO 104 (to be compensated) to realize temperature compensation.\nWhere the temperature-compensated voltmeter is constructed in advance by collecting the voltage of the VCXO 104 at different temperatures while maintaining the frequency stabilization. A more detailed description of the construction can be found in following references:    1. Huang X, Liu D, Wang Y, et al. 100-MHz Low-Phase-Noise Microprocessor Temperature-Compensated Crystal Oscillator[J]. Circuits & Systems II Express Briefs IEEE Transactions on, 2015, 62(7):636-640;    2. “Temperature compensation for an oscillator crystal”, Inventor: Markus Hammes etc, US Patent Publication Number: US 20170085271A1, Date of Publication: Mar. 23, 2017;    3. “Digitally compensated phase locked oscillator”, Inventor: Nicholls Charles William tremlett, etc, US Patent Publication Number: US 20160365865A1, Date of publication: Dec. 15, 2016.\nIn sum, the TCXO in prior art uses an open-loop compensation architecture, a temperature sensor is needed, and the temperature sensor should be as closer as possible to the crystal resonator on a circuit. However, the resonant wafer of the crystal oscillator is individually enclosed in a confined space, which inevitably produces a temperature hysteresis between the temperature sensor and the resonant wafer, leading that there is no significant breakthrough in the TCXO frequency temperature characteristics. Especially for the crystal oscillator with high frequency output, the temperature hysteresis is more obvious."}
{"text": "Typically, in bathrooms, mirrors are either rigidly mounted to the wall surface or contained as part of a vanity such that the mirror is fixed in a generally vertical orientation with respect to the vertical axis of the bathroom wall surface. However, with the advent of federal, state and local regulations regarding rights of the handicapped, such as the American with Disabilities Act, \"accessibility\" of bathroom fixtures for use by the handicapped is of great importance, especially in public settings. In lieu thereof, it has been found desirable to provide a mechanism which allows the bathroom mirror to be easily tilted at a predetermined angle generally downwardly with respect to the vertical axis of the wall surface such that it is accessible for use by handicapped individuals.\nIn the past, in order to provide a pivotally tiltable mirror, typically, the mirror assembly has included two knobs extending outwardly from the side edges thereof. These knobs fit into generally U-shaped slots provided at the ends of upwardly extending arms of a support assembly. In order to pivotally tilt these mirrors, either the upper or lower portion of the mirror is pushed forwardly or rearwardly to obtain the desired mirror angle. These mirrors have been found disadvantageous, especially in public places, as the pivotal angular movement of the mirror is not restricted thereby leading to potential breakage of the mirror due to accidental contact of the mirror assembly with another surface as well as inconvenience in terms of the user maintaining the mirror assembly at the desired angle. Moreover, these mirror assemblies are not practical for use in bathroom settings. It has therefore been found desirable to provide a tilting mechanism for bathroom mirrors which allow the bottom portion of the mirror to be pushed inwardly such that pivotal angular movement of the mirror is restricted to thereby minimize the possible attendant chances of mirror breakage while at the same time retaining the mirror in the proper angular position."}
{"text": "Ethanol for industrial use is conventionally produced from petrochemical feed stocks, such as oil, natural gas, or coal, from feed stock intermediates, such as syngas, or from starchy materials or cellulose materials, such as corn or sugar cane. Conventional methods for producing ethanol from petrochemical feed stocks, as well as from cellulose materials, include the acid-catalyzed hydration of ethylene, methanol homologation, direct alcohol synthesis, and Fischer-Tropsch synthesis. Instability in petrochemical feed stock prices contributes to fluctuations in the cost of conventionally produced ethanol, making the need for alternative sources of ethanol production all the greater when feed stock prices rise. Starchy materials, as well as cellulose material, are converted to ethanol by fermentation. However, fermentation is typically used for consumer production of ethanol for fuels or consumption. In addition, fermentation of starchy or cellulose materials competes with food sources and places restraints on the amount of ethanol that can be produced for industrial use.\nEthanol production via the reduction of alkanoic acids and/or other carbonyl group-containing compounds has been widely studied, and a variety of combinations of catalysts, supports, and operating conditions have been mentioned in the literature. During the reduction of alkanoic acid, e.g., acetic acid, other compounds are formed with ethanol or are formed in side reactions. These impurities limit the production and recovery of ethanol from such reaction mixtures. For example, during hydrogenation, esters are produced that together with ethanol and/or water form azeotropes, which are difficult to separate. In addition when conversion is incomplete, unreacted acid remains in the crude ethanol product, which must be removed to recover ethanol.\nTherefore, a need remains for improving the recovery of ethanol from a crude product obtained by reducing alkanoic acids, such as acetic acid, and/or other carbonyl group-containing compounds."}
{"text": "1. Field of the Invention\nThe present invention relates generally to food processing apparatus and, more particularly, but not by way of limitation, to a food processing apparatus for removing leaf fat from the abdominal cavity of a hog carcass.\n2. Description of the Prior Art\nIn the past, various devices have been constructed for removing various parts from various types of animal carcasses and the like.\nThe U.S. Pat. No. 3,714,681, issued to Van Snellenberg disclosed a device for cutting glands from membranes of fish. The gland was pulled away from the membrane by a suctionhead and a rotating cutter sleeve rotating about the sectionhead was automatically moved into position for cutting the membrane from the gland. The U.S. Pat. No 1,997,843, issued to Warrell, disclosed a device for producing openings in multi-ply flexible materials wherein a tubular cutter was automatically moved into rotating engagement with the material. The British Pat. No. 833,422, issued to Burgess disclosed a device for boring encysted areas from fillets of fish.\nVarious other devices for processing animal carcasses and the like were disclosed in the U.S. Pat. Nos. 3,849,836, issued to Bernard; U.S. Pat. No. 3,820,196, issued to Penman; U.S. Pat. No. 3,685,095, issued to Metro; U.S. Pat. No. 3,533,131, issued to Ivarsson; U.S. Pat. No. 3,234,591, issued to Vogt et al.; U.S. Pat. No. 3,177,520, issued to Vogt et al.; U.S. Pat. No. 3,156,945, issued to Vogt et al.; U.S. Pat. No. 3,129,456, issued to Renfroe; U.S. Pat. No. 692,824, issued to Burleson; and German Pat. No. 237,892."}
{"text": "This invention relates to an improved transducer for use in a vibratory viscometer; which transducer is particularly useful for, but not limited to, the inline detection under flow conditions of viscous and elastic properties of fluids being processed.\nVibratory viscometers are well known in the art and generally employ a transducer which has an immersible portion which is vibrated with a small amplitude. Fluid viscosity/density/viscoelasticity can be determined from the frequency and/or amplitude changes in the vibration and/or the power required to sustain the vibration when the immersible portion of the transducer is immersed in a fluid.\nSuch transducers generally comprise (i) an immersible tip, (ii) an electromagnetic drive and (iii) an electromagnetic or piezoelectric pickup. Transducers of this type are described by J. G. Woodward in \"Vibrating Plate Viscometers\", Electronics, February 1952; and by J. D. Ferry in \"Viscoelastic Properties of Polymers\", published by John Wiley & Sons, New York, 1970.\nPrior Oscillatory Viscometers\nThe oscillatory viscometers described below exemplify the prior art, and have transducers which interact with the fluid being measured by \"surface loading\" of an oscillating portion of the transducer; said term having been used J. D. Ferry in his aforementioned publication.\nThe oscillating surface generates shear waves in the liquid or other fluid. Viscoelasticity is measured in terms of the characteristic impedance of the liquid.\nIn the vibrating plate viscometer described in the aforementioned publication of J. G. Woodward, the moving end of a vibrating reed supports a plate which detects viscous resistance of the fluid in which it is immersed. The reed is driven electromagnetically and the oscillations are picked up by a barium titanate piezoelectric block. The reed is clamped at one end, exposing the driver and pickup to fluid vapors. The viscometer measurements are confined to viscous loss determined from the decrease in amplitude of oscillation observed when the plate is immersed. The plate vibrates in a direction essentially perpendicular to its major surfaces.\nA. Konno, S. Malsino and M. Kameko, in Japan Journal of Applied Physics 189 (1968), reported their measurements of storage modulus and viscous loss by oscillating an immersed very thin microscope slide at 100 Hz. This apparatus, driven by a moving coil, is reported as strictly a research instrument. The internal workings of the transducer were exposed to fluid vapors.\nA viscometer known as Le Viscosimetre \"Pivert\" is sold by the Societe Francis de Service, 8 rue Nobel Zl, 45700 Villemander, France. The sensor tip of this viscometer (also called \"Sofraser\") is a U-shaped stainless steel needle. One leg of the U-shaped needle is mechanically driven sinusoidally at 125 Hz., which is near the mechanical resonance frequency of the transducer. Both legs of the needle are clamped at nodal points. Frequencies are measured at both legs. The phase difference between the frequencies is transformed into a voltage or current related to the viscosity of the liquid in contact with the needle. The manufacturer states that installation conditions can decrease accuracy and therefore should be carefully checked to be sure of optimum accuracy. The range of viscosity measurement is 0.5 to 15,000 mPa.s (cPs). The environmental limits of utilization are 200.degree. C. and pressure to 100 bars.\nThe Dynatrol Viscosity Detector is sold by Automation Products, Inc., 3030 Max Roy, Houston, Tex. 77008. Although this instrument does not employ a blade sensor, it does employ welding at the node point of a resonating system in order to isolate the oscillating probe from the driver and pickup. Viscosity-density product is detected by an immersed 5 inch long stainless steel probe. The probe is made from a stainless steel rod and bent very much like a hairpin. Both ends of the rod are welded at nodal points to a supporting plug. One end of the probe penetrates into the housing where it is electromagnetically driven to resonate in flexure at 120 Hz. The other end of the probe also penetrates through the plug into the housing where there is an electromagnetic pickup coil. The decrease in amplitude of vibration of the immersed probe due to interaction with the fluid is electronically converted to viscosity-density product.\nThe Labor-Viskosimeter QV35 is marketed by Bopp & Reuther GmbH, Car-Reuther-Strasse 1, Postfach 310140, D-6800 Mannheim 31, West Germany. This instrument uses an oscillating transducer to measure viscosities of laboratory samples of liquids by means of a quartz crystal sensor oscillating in torsion at 55 KHz. The damping effect of the fluid being measured on the amplitude of oscillation is converted into viscosity. The temperature range of this instrument is limited to -50.degree. to 150.degree. C.\nThe Model 1800 Viscometer sold by Combustion Engineering, Inc., P.O. Box 831, Lewisburg, W.Va. 24901 utilizes an oscillating sensor blade. Short pulses of current are applied to the top of the blade, which is composed of a magnetostrictive alloy. The blade protrudes through a metal diaphragm. Each pulse causes the blade to vibrate at its natural frequency of 28 KHz. When the amplitude of vibration of the immersed blade has fallen to a preset value that relates to the viscosity of the liquid, another pulse is automatically applied. The change in pulse rate is proportional to the square root of viscosity-density product. The range of viscosities measured is from 0 to 5,000 centipoise x grams/cm..sup.2 in ranges of 0-50, 0-500, and 0-5000. This viscometer was produced for many years by Bendix Corporation, and was first described as the Ultra-Viscoson by W. Roth and S. R. Rich in Jr. Applied Phy. 24 940-950, July 1953. The Bendix Ultra-Viscoson is described on page 308 of Viscosity and Flow Measurement by S. R. VanWaser, S. W. Lyons, K. Y. Kim and R. E. Cowell, Interscience Publishers, New York, 1963. Shortcomings included very high frequency of measurement, fragility of very thin strip (blade), lack of sensitivity at very low viscosities, need to flick the strap from time to time to relieve strains in the magnetostrictive alloy; and since strips were easily bent, replacement blades needed to be available.\nHermetic sealing between driver and pickup is a feature of the \"Vibrating Sphere\" and \"Viscoliner\" oscillating viscometers of Nametre Company, 101 Forrest St., Metuchen, N.J. 08840, the assignee of the present application. However, the spherical and cylindrical sensors are somewhat obstructive to flowing fluids, particularly slurries. In order to reduce the obstruction to flow, the diameter of the pipe in which the fluid flows needs to be increased so as to satisfactorily accept the oscillating sensor. Since the mode of oscillation has two degrees of freedom, care must be taken to chose between in-phase and out-of-phase torsional motion. Further, the oscillating surface generates diverging shear waves that may be so long in wavelength that they are not conveniently reflected for accurate measurement of highly viscoelastic fluids.\nU.S. Pat. No. 4,729,237 describes a tuning fork transducer, each of the two arms of the fork having welded to it a blade that is oscillated in the liquid. It is claimed that one blade on one arm of the tuning fork gave less accurate viscosity measurements than having blades on each arm. This viscometer is a laboratory instrument. The specification describes vertical motion of a sample container to immerse the two blades in the liquid.\nObjects of the Present Invention\nFor a long time there has been great need for small inline blade sensors that do not impede flow and that are intrinsically separated from the fluid being measured. Examples of these needs include transport of mineral slurries such as powdered coal and powdered lime where particle concentration must be controlled. The consistency of food fluids such as ice cream, coffee and bread dough needs to be monitored and controlled by means of reliable rugged viscosity sensors. Better control of viscoelastic properties of polymer fluids can be achieved by employing suitable sensors. Accordingly, an object of the present invention is to provide an inline blade sensor for a vibratory viscometer that impedes fluid flow to a substantially lesser extent than prior art sensors. Another object of the invention is to provide such a sensor in a configuration that facilitates isolation of the drive and pickup portions of the transducer from the fluid.\nOther objects of the invention are to provide transducer-sensor devices that are relatively small, are rugged, easily inserted into industrial process pipes and equipment; that operate over a broad range of viscosities under hot and cold conditions and at high pressures; that can be speedily installed and quickly operated in remote and dangerous locations; that are dependable over long stretches of time; that can be easily removed, cleaned and reinstalled; that are capable of providing viscous loss and/or transducer sensor frequency signals for conversion into loss modulus and storage modulus values; that are capable of being vibrated over ranges of frequencies as well as at mechanical resonance; and that are adaptable for use in very viscous and very elastic liquids."}
{"text": "Vertical power transistors, such as group IV based trench type field-effect transistors (trench FETS), are used in a variety of applications. For example, silicon based trench metal-oxide-semiconductor FETS (trench MOSFETs) may be used to implement a power converter, such as a synchronous rectifier, or a direct current (DC) to DC power converter.\nFor many trench MOSFET applications in which high switching speeds are necessary or desirable, it is advantageous to reduce gate charge so as to reduce switching loss. In addition, as the electronic devices and systems utilizing power MOSFETs are continually reduced in size, there is a corresponding need to reduce the size of the MOSFETs themselves. However, the resulting reduction in cell pitch and increased power density tends to undesirably increase gate charge in conventional trench MOSFET structures."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates to battery housing and, particularly, to a battery cover latching assembly used in a portable electronic device.\n2. Description of Related Art\nBattery cover latching assemblies are usually provided to secure and allow required access to batteries to portable electronic devices.\nA typical battery cover latching assembly for an electronic device usually includes a pair of hooks on one end of a battery cover of the electronic device and a locking pin on the other end of the battery cover. Correspondingly, a pair of grooves is defined in a housing of the electronic device. A locking hole corresponding to the locking pin is also defined in the housing. By attaching the battery cover onto the housing, the hooks can enter into the grooves, and the locking pin is securely latched into the locking hole. Thus, the battery cover is secured to the housing.\nHowever, to detach the battery cover, considerable force is needed due to the tight latching of the locking pin into the locking hole. Thus, the battery cover may be damaged, and it is difficult to replace the locked battery.\nTherefore, there is room for improvement within the art."}
{"text": "Recently, there has been a resurgence of interest in artificial neural networks (ANN), and such research has commonly been termed deep learning, computer learning and similar terms. The increase in general-purpose processor computation power has given rise to the renewed interest that waned a couple of decades ago. Recent applications of ANNs have included speech and image recognition, along with others. There appears to be an increasing demand for improved performance and efficiency of computations associated with ANNs."}
{"text": "There are many applications which require limiting voltage levels of alternating or oscillating voltage sources.\nFIG. 1 shows a simplified circuit diagram of a limiter according to the prior art. There is a voltage source VS represented by an ideal voltage source VC and an impedance RI. The voltage source VS may be, e.g. an oscillator, and most specifically, an oscillator in an RFID application. The output of the voltage source VS is applied to electronic components. A first transistor P1 is coupled with its channel between the output node OUT and ground. A second transistor N1 is coupled with its channel between the output node OUT and a supply voltage level VDD. The gates of the first transistor and the second transistor N1 are biased by BIAS voltage levels VBP1 and VBN1 respectively. However, any currents drawn from the output node OUT through transistors N1 or P1 are finally dissipated and contribute to the overall power consumption."}
{"text": "The invention relates to an application of mechanics to medicine. It deals with improvements in vacuum dosing devices used in surgical centers and other hospitalization areas including intensive care units of hospitals.\nThe invention introduces a series of improvements in the operation of the instrument described in U.S. Pat. No. 13,437 (Ser. No. 23,472). Although the instrument disclosed in the above mentioned Patent has a vacuum saving function, it shows some imperfections in its operation as noted below.\n(a) The closure produced by flattening the rubber tube is not perfect. As it collapses, said tube leaves in both corners (FIG. 1xe2x80x941), a space along the tube, through which the vacuum unit continues drawing in ambient air during the stand-by state. Due to its flattened shape, the rubber tube frequently does not perfectly fit in the tube cylinder thus leaving a space (2) which continues with the space mentioned in (a), and through which air is permanently lost.\nThe hole (5) that communicates the ambient air with the inside of the case (3), through the hole (4) in the case, and the latter with the vacuum system, permits a vacuum loss which, though small, is permanent.\nAs the aspirated material passes in front of the hole (5), a certain amount of said material (6) enters the case (3) contaminating the finger that occludes its hole.\nDuring the production process, both solid pieces need to be connected to a soft rubber tube, and subsequently the resulting set, which is also soft, must be inserted and fixed to the case holes; this assembly operation is therefore complicated and slow."}
{"text": "Exposure apparatuses are commonly used to transfer images from a reticle onto a substrate during the manufacturing and processing of liquid crystal displays (“LCDs”) and semiconductor wafers. There is a never ending desire to manufacture larger LCDs. Typically, larger masks are required to manufacture larger LCDs. Unfortunately, as the size of the masks utilized increases, so does the likelihood that the mask may be subject to a certain amount of sagging due to gravity in the middle region of the mask that is not directly supported. Accordingly, there is a need to develop a system whereby the potential sagging of the mask is minimized while inhibiting unwanted blocking of the pattern on the mask."}
{"text": "In a rotating disk computer mass storage device, data and control information is typically recorded on concentric tracks on the disk. One type of arrangement imbeds the control information in servo fields. The servo fields are spaced periodically along the length of each track. The servo fields include information which describes the track number, the position on the track, synchronization information, and other information. Between each of the servo fields, data is recorded in a data field. Each servo field is utilized to identify the following data field and to position accurately a head or transducer which reads or writes the data field, among other functions. This type of arrangement is referred to as an imbedded servo field and control system because the servo field is imbedded in each of the tracks, rather than other arrangements which record the servo and control information on separate tracks from those which contain data.\nTypically, the information contained in the servo fields is recorded by a precise laser/magnetic recording or other suitable device at the time of manufacture of the disk drive. By using the servo field information, the user may then write data into the data fields interspersed between the servo fields, or may access and read data previously recorded between the intervening servo fields.\nProper functionality of a rotating disk computer mass storage device requires that the disk be rotated at a substantially constant speed. Servo system controllers are employed to, among other functions, regulate the rotational speed of the disk. However, due to the practical response limits of a control system, it is sometimes very difficult or impossible to maintain a high degree of precision in controlling the rate of disk rotation. Furthermore, because the frequency at which information is recorded on each track is extremely high compared to the rotational speed of the disk, slight variations in disk rotational speed translate into significant variations in the signalling frequency for writing to or reading from each track. These signalling frequency variations present significant difficulties when attempting to synchronize the logic of the servo system controller to the asynchronous information received from the disk. Thus, despite very effective motor speed regulation control systems, it is virtually impossible to avoid signalling frequency variations.\nThe problems of controlling the rotational speed and synchronizing the reading of data from or writing data to rotating disk mass storage devices can be particularly difficult in portable computers. Portable computers are frequently used on unstable surfaces, and therefore are subject to being physically disturbed and moved. Each physical movement of the computer applies forces on the rotating disk, and those forces typically cause minor perturbations or variances in the speed of the disk.\nEven in circumstances where physical movement of the computer does not translate into variations in disk rotational speed, other factors may intervene. For example, despite the normally reliable and consistent manner in which the servo field information is recorded on the tracks, occasionally there are variations in the frequency or periodicity at which the information is recorded. These variations may cause variations in timing and synchronization even though the computer disk is rotating at precisely the desired rate.\nThe information recorded on the tracks is in the form of a curvilinear series of dibits, or magnetic field reversals. The strength of the magnetic field reversals is very small, thereby presenting a situation where electrical noise may represent a significant component of the signals induced in a transducer or head of the disk drive. The presence of noise raises the risk that the noise may obscure or corrupt the recorded information. The noise may cause substantial difficulty in correctly detecting and responding to the servo information when controlling the operation of the disk drive. In some situations, noise may be the most or a highly significant detriment to the best operation of the servo system controller and the disk drive device."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of brooms, and in particular push brooms and handheld whisk brooms.\n2. Description of the Prior Art\nBrooms in general are old and well known in the art. Push brooms are known and handheld whisk brooms are known. However, no embodiment known to the present inventor incorporates the concept of a convertible broom where a push broom and a whisk broom are incorporated into the same device."}
{"text": "1. Field of the Invention\nThe present invention relates to a system manufacturing composite aluminum panels continuously by using metal (aluminum) cores being expanded continuously in manufacture of composite aluminum panels.\n2. Description of the Related Art\nUp to now, for manufacture composite aluminum panels inserted honeycomb cores, honeycomb type metal (hereinafter refer to as aluminum) cores are expanded in a certain length by a manual expanding device, hot-melt film layers are put on the upper and the lower sides of the expanded cores and on the film layers aluminum plates are put, the objects are pressed in hot presses, to give composite aluminum panels.\nBut in this case expanding devices of aluminum cores should be operated manually and thus the cores can not expanded continuously.\nAccordingly the standardized panels can not be manufactured continuously.\nNow, referring to the accompanying drawings, a method for expanding cores for composite aluminum panel by the conventional device for expanding cores and a conventional process for preparing panels by using the expanded cores will be described as follows.\nAs shown in FIG. 5, a fixing part Fa fixing a left end of pre-expanded raw material 1001 for honeycomb type cores comprising a supporter 300 fixed on the upper one side of a base 200, a number of air pressure cylinders 400 ascending and descending on the supporter 300, a stair type bracket 600 fixed on the lower ends of rods 500 of the cylinders 400, a holding rod 900 for taking needle rods 700 mounted on the bracket 600 and having a number of holes 800 capable of inserting and drawing needle rods 700 for fixing on a number of holes 1002 of honeycomb type cores; a transferring part R comprising the supporter 300, a number of the air pressure cylinder 400, the stair type bracket 600 and a holding rod 900 for taking needle rods having a number of inserting and drawing holes 800 identical to those of the fixing part Fa, and a pinion gear (not shown) established on the lower end of the supporter 300, engaging with a rack 1003 established on the upper both sides of the base 200 to transfer right and left; and plate 100 for expanding the honeycomb type core formed between the fixing part Fa and transferring part R.\nBut to manufacture the expanded honeycomb type cores from raw material for core by the conventional device for expanding honeycomb type core, after a number of needle rods 700 are inserted perpendicularly in the inserting and drawing holes 800 for the needle rods of the fixing part Fa manually, by means of ascent and descent of ascending and descending cylinder the bracket 600 fixed on the lower part of the cylinder is gone up and come down to ascend and descend the holding rod 900 for taking needle rods 700 fixed on the bracket 600, and by means of ascent and descent of the holding rod 900 the lower ends of the needle rods 700 inserted in the inserting and drawing holes 800 of the holding rod 900 are fixed, by means of inserting, in the holes 1002 of the honeycomb type core expanded partially in left end of raw material for the honeycomb type core.\nAt this time, by means of the pinion gear (not shown) fixed on the lower end of the transferring part R engaged with the rack established on the upper both sides of the base 200, the transferring part R is transferring manually to lie adjacent to the fixing part Fa, then the lower ends of the needle rods 700 inserted in the inserting and drawing holes 800 for needle rods of the holding rod 900 fixed on the transferring part R are inserted in the holes 1002 of the honeycomb type core expanded partially or manually of the right end of the raw material (core-raw material) 1001 for the honeycomb type core and the transferring part R is transferred manually in the right direction to give expanded honeycomb type core from the raw material 1001 for the honeycomb type core.\nHowever, in the conventional device, the needle rods 700 for taking the core-raw material 1001 should be inserted manually in the inserting and drawing holes 800 for needle rods, the transferring part R also should be transferred manually in right and left, length and width of the honeycomb type core manufactured by the device also are restricted, that is, in one device cores having restricted length and width only are manufactured.\nAccordingly this device has faults that various sorts of cores having different length and width cannot be made.\nFurthermore, in manufacturing composite aluminum panels by using honeycomb type cores expanded by this device because all of honeycomb type cores, hot melting films and aluminum plates as raw materials for manufacturing the panel should be layered only manually, the raw materials cannot be layered in the same mode.\nAccordingly by means of the above-mentioned conventional device standardized articles cannot be produced and also long and large panels cannot be produced continuously."}
{"text": "This invention relates to a remote control apparatus and, more particularly, to an axial alignment aid and related method for facilitating remote control operations.\nMany attempts have been made to remotely align two objects in space. For example, the National Aeronautics and Space Administration (NASA) has done extensive work in the area of docking spacecraft. As discussed in \"Communications, Tracking and Docking on the Space Station,\" by Erwin et al., published by the Institute of Electrical and Electronics Engineers (1982), NASA has proposed the use of a triangular array of reflectors on a passive target to define a plane and facilitate docking of an active vehicle thereto. Sequential tracking of each reflector from the active vehicle in space allows altitude determination between the active vehicle and the passive target. However, this triangular reflector system necessarily employs computers, thus making it extremely complex. In addition, the system is incapable of accurately and remotely aligning an axis of importance of the passive target to an axis of importance of the active vehicle to allow \"soft\" docking.\nOther attempts have been made to accurately align or dock one object relative to another. For example, U.S. Pat. No. 3,269,254, issued to COPPER et al., discloses an optical apparatus including: a light source; a beam splitting cube; a lens system and a grated mirror, all on one side of the light source; and a measuring reticle having coordinates thereon located on the opposite side of the light source. Light from the light source is directed by the beam splitting cube through the lens system. The light then reflects off of the mirror to again pass through the lens system and the beam splitting cube. Finally, an image is formed at the measuring reticle to provide information regarding orientation of a body o which the mirror is attached, as defined by azimuth, pitch and roll. Because the mirror described in this patent is flat, if the beam of light directed at the mirror is too far off-axis, the light beam can be reflected too far off, thus preventing the formation of an image at the measuring reticle. As a result, an accurate off-axis measure is not available with this patented device. Further, this patented device is sensitive to distance, i.e., the effectiveness of the device is lessened as the measuring reticle and the mirror are placed farther apart. Finally, this device is not adapted to allow docking of one member relative to another member.\nIn addition, U.S. Pat. No. 2,352,179, issued to BOSLEY, uses a set of four photo-electric cells as detectors to provide information as to the orientation of a lens. That is, the set of cells indicates from one perspective how far away one is from a target; but such detectors cannot accurately indicate the axis that is defined by the target. Therefore, such a system using a set of four detectors is also incapable of allowing on-axis docking alignment of one object relative to another.\nOne application where remote control docking would be particularly important is in nuclear powered steam generators. For example, performing a task within the steam generator with a tool or end effector and a robot for receiving the end effector is currently performed as follows. A first set of cameras positioned in the steam generator allow a conventional robot arm known as ROSA (remote operation service arm) to be remotely viewed and moved to the general vicinity of one or several end effectors positioned in the steam generator. These cameras, however, cannot accurately and automatically align the respective axes of the chosen end effector and ROSA, i.e., cannot reliably effect docking and attachment.\nThe end effector also has a camera mounted thereon which ultimately is used to remotely view the worksite during performance of the task via a conventional feedback system using a remote viewing station, robotic controls and means for moving ROSA.\nAs a result, a man must be introduced into the steam generator to manually align and attach the end effector to ROSA. However, due to the potential radiation exposure, tedious alignment and attachment steps necessary, and the relatively heavy members being manipulated, it is desired to make the alignment and attachment of end effectors to ROSA in the steam generator an entirely remotely controlled operation having high reliability A more general need also exists for a means by which accurate alignment can be made between two important axes of respective docking members using single perspective visual guidance."}
{"text": "The present invention is directed to an electronic digital readout diet scale for measuring and displaying weight information.\nPresent bathroom and other scales generally are a very poor aid to weight control because they do not provide any psychological reinforcement to dieters because the scales do not indicate the progress being made after several days of dieting. In other words, these scales are unable to detect and indicate progress made as a result of a diet program.\nA prior art diet scale is shown in U.S. Pat. No. 3,967,690 which issued on July 6, 1976 to the inventor of the present invention. The subject matter of that patent is hereby incorporated by reference. U.S. Pat. No. 3,967,690 discloses a digital readout diet scale which accurately measures and displays changes in the weight of a dieter. The present or current weight of a dieter is measured and this weight is then compared to a previous weight in order to determine the change in weight. A circuit arrangement is disclosed for accomplishing these weight measurements.\nThe present invention is an improvement of the digital readout diet scale of the type disclosed in the above United States patent. In the present invention, all the hardwired circuitry of the above patent has been eliminated. This invention also contemplates several additional functions. In particular, as set forth below, a microcomputer is used to perform the various functions previously performed by the discrete components of the above patent. In addition, the diet scale of the present invention enables a dieter to enter a diet program in the microcomputer having an objective weight and given time interval and the diet scale will indicate the amount of weight to be lost or gained in order to achieve the objective weight as well as the deviation between the current weight and the projected weight according to the diet program."}
{"text": "(a) Technical Field\nThe present invention relates to a control method for improving fuel efficiency of a hybrid electric vehicle. More particularly, the present invention relates to a method for controlling a fuel cut operation to be performed in a wider transmission region by a reaction control of a motor, thus improving fuel efficiency of a hybrid electric vehicle.\n(b) Background Art\nIntensive research has been underway to provide a high efficiency hybrid electric vehicle (HEV) by improving fuel efficiency and emission performance.\nOne of the various techniques for improving fuel efficiency and emission performance is to provide an HEV with idle stop system, in which the engine of the vehicle is stopped while the engine is idling as the vehicle is stopped.\nAnother technique for improving fuel efficiency and emission performance is to increase fuel cut operation. A fuel cut mode saves energy by stopping fuel injection when a vehicle running at a high speed is decelerated.\nHowever, it has a problem of the fuel cut range being restricted. More specifically, a sudden change in engine torque can occur when the engine is restarted as the fuel is injected. When a torque is applied at a low speed, for example, at a low gear 1 or 2, the vehicle may be shaken due to the sudden change in engine torque.\nFor example, when the engine speed is lowered below about 1500 RPM, a fuel in (fuel re-injection) is performed. Although it is desirable to increase the fuel cut period for improvement of fuel efficiency, impact is brought about at the moment when the fuel cut operation is cancelled (i.e., rewetting is started). To improve fuel efficiency, it is necessary to prolong the fuel cut period at a lower speed and RPM.\nIn the event that a fuel cut operation is performed for an ordinary vehicle, the fuel cut operation is cancelled at about 1500 RPM to revive the engine. Converting the RPM into a vehicle speed, the fuel cut operation is cancelled at about 60 km/h (and 1500 RPM) in case of a vehicle not exceeding 2000 cc, although it may vary according to a gear ratio of the vehicle.\nAccordingly, if the vehicle is running at 60 km/h or higher for a predetermined period of time, the fuel cut effecting condition is permitted. The fuel cut operation continues during deceleration period and is cancelled when a preset condition is met.\nIn detail, the fuel cut cancellation is made when it reaches a predetermined RPM with respect to coolant temperature, and a corrected RPM is considered thereto. Moreover, in the case where an air conditioner is turned on, an additional RPM is considered and, in the case where the air conditioner is turned off, the fuel cut cancellation is made at about 60 km/h and 1500 RPM.\nHEVs developed so far reach about 1500 RPM at 60 km/h and the vehicle speed is decreased while maintaining about 1200 RPM up to about 15 km/h. With the current fuel cut logic for gasoline vehicles, it is thus impossible to improve fuel efficiency of HEVs. Moreover, a torque of about 15 to 25 Nm is suddenly applied during the fuel cut cancellation, which results in a serious impact that a driver can feel when the gear ratio is low.\nThe information disclosed in this Background section is only for enhancement of understanding of the background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art that is already known to a person skilled in the art."}
{"text": "This invention in general relates to the void fraction measurement in gas/liquid two-phase flow system. Particularly, the invention relates to a method and a pressure-differential device to measure the void fraction of a gas/liquid two-phase stratified flow within a horizontal pipeline system. The said method and said device are applicable to flow at steady and transient states, and at high-temperature/pressure conditions.\nVoid fraction is a crucial parameter in the determination of energy and momentum transfer in a two-phase flow system. In usual practice, the void fraction is defined two-dimensionally, in reference to a chosen cross section of the pipe under measurement, as the area of gas phase passing through that cross section divided by the whole cross-sectional area. At present, the instruments most often used for void fraction measurement are the expensive and complicated xcex3-ray and \"khgr\"-ray densitometers, in which void fraction is derived from radiation attenuation of one or several xcex3-ray or \"khgr\"-ray beams traversing through the gas/liquid two-phase flow. FIG. 6 illustrates such a densitometer commonly used in industry, which includes radiation source 60, collimators 61, two-phase flow measurement section 62, radiation shields 63, and a detector and counting system 64. This layout was first reported by Edelman et. al. in International Journal of Heat and Mass Transfer, vol. 28, no. 7, 1985, where xcex3-ray is applied to measure the average cross-sectional void fraction of a steady-state boiling flow within a stainless steel pipe. FIG. 7 is another example, presented by Kagawa et. al. in International Topical Meeting on Nuclear Reactor Thermal Hydraulics, NUREG/CP-0014, vol. 2, 1990, where \"khgr\"-ray scanning method is applied to a high-temperature/pressure flow system to measure the distribution of cross-sectional void fractions in a horizontal pipe and the system is at transient state due to a severe leakage somewhere in the pipeline system. This is an even more complicated instrument than the xcex3-ray densitometer just mentioned, in that, more than those basic component parts shown in FIG. 6 the instrument has, referring to FIG. 7, a radiation scanning location-control mechanism 65 and a radiation scanning beam-control electronic system 66. The measurement accuracy of xcex3-ray or \"khgr\"-ray densitometer depends on the degree of homogeneity of the gas/liquid two-phase flow: better homogeneity of the fluid obviously promises more reliable measurement result. However, it is well known that, the void fraction of a gas/liquid two-phase flow seldom is, if ever, homogeneously distributed. The measurement accuracy of a radiation densitometer also depends on the incident angle and location of each radiation beam, because each beam, different from the rest, is specifically attenuated by a certain wall thickness of pipe and by a certain portion of gas/liquid two-phase flow. As a rule, the farther away is a radiation beam from pipe axis, the less reliable is the measurement data. Moreover, depending on the material and thickness of pipe wall, the radiation attenuation may be so great as not to leave any useful information for the detector to resolve. For a specific example, one may consider the case of a horizontal natural circulation of a gas/liquid two-phase thermal-hydraulics system, in which the density difference between low flow rate liquid phase (e.g. water) and high flow rate gas phase (e.g. vapor) is usually so large that, under the influence of gravity, the two phases virtually separate and become two quite distinct flows, namely the upper layer gas flow and the lower layer liquid flow, thus the term flow stratification. Neither xcex3-ray nor \"khgr\"-ray densitometer is at all suitable for such flow condition for the want of fluid homogeneity.\nTo overcome the aforesaid deficiencies of radiation method for void fraction measurement, the present invention, taking advantage of the fact that flow stratification phenomenon often occurs in natural circulation of gas/liquid two-phase flow within horizontal pipeline system, discloses a method to determine the void fraction of gas/liquid stratified flow; wherein a couple of pressure-signal tubes are vertically flush-mounted to the top and bottom interior surface of the horizontal pipeline to measure the pressure differential across the pipe. A thermocouple-rake temperature probe assembly is used to obtain the average temperature of the liquid phase of a gas/liquid two-phase stratified flow, in which the thermal stratification also occurs. Therefore, the average mass density of the liquid phase of the stratified flow is obtained. Then the liquid-phase level of the stratified flow is calculated from the aforesaid pressure-differential data and liquid-phase mass density data. Finally the void fraction is derived from the known geometric relation between said liquid-phase level and the cross section of the pipe.\nSince commercialized instruments for void fraction measurement of gas/liquid two-phase flow, such as xcex3-ray or \"khgr\"-ray densitometer, are far from satisfactory when flow stratification occurs, owing to the reasons briefly discussed in the foregoing, it is the principal objective of this disclosure to resolve this problem by providing a novel method for void fraction measurement, applicable to vapor/water, or any other kind of gas/liquid two-phase stratified flow within a horizontal pipeline system. According to this disclosure, a couple of pressure-signal tubes are vertically flush-mounted to the top and bottom interior surface of the pipe at a selected location, so as to measure the pressure differential across the cross section at that location; and, through liquid-phase mass density correction, this pressure-drop data is transformed into the liquid-phase level within the pipe; and the said liquid-phase level, by its geometric relation with the cross section of the pipe, is then transformed into void fraction. The device herein disclosed can accurately determine the void fraction of a horizontal gas/liquid two-phase stratified flow, be the flow at steady or transient state, and whether or not the flow is at high-temperature/pressure condition.\nWhen the flow being measured is at high-temperature/pressure condition or is at transient state, physical effects like gravity and heat that may render pressure-differential measurement difficult must be put into consideration. Not only does gravity produce flow stratification, but it may also pull down some pressure-signal transmitting medium from the upper pressure hole, resulting in voids or even void sections present within the upper pressure-signal tube. To obviate the gravitational effect four measures are taken: (1) keep as small as possible the diameter of pressure-signal tube that is near the pressure hole, (2) keep as short as possible the vertical section of the upper pressure-signal tube that is near the pressure hole; (3) keep as small as possible the two locations at top and bottom pressure holes where pressure-signal transmitting medium is present; and (4) have pressure-signal tubes flush-mounted to the top and bottom interior surface of the pipe. As to heat, obviously it can increase the temperature of pressure-signal tubes, especially at the sections near the pressure holes, because of thermal conduction. A density gradient will therefore occur within pressure-signal transmitting medium. This problem will be especially serious at the upper vertical section of the pressure-signal tube, and will greatly affect measurement accuracy of the liquid-phase level of a gas/liquid two-phase stratified flow within a horizontal pipeline system. If accident such as leakage or cleavage occurs in a high-temperature/pressure closed system, resulting in pressure decrease faster than temperature decrease, vaporization of the liquid phase within the system may begin rapidly (this is called flashing phenomenon), and this may sometimes lead to vaporization of pressure-signal transmitting medium near the pressure holes. Flashing of pressure-signal transmitting medium may also result in voids or even void sections within pressure-signal tube. Once voids or void sections are present within the vertical section of pressure-signal tubes, any further measurement is not merely futile but, in some cases, possibly dangerous may result. This is because the measurement device is no longer reliable: it may produce false signals which may accordingly trigger the safety systems of, say, a nuclear reactor to work, or, contrariwise, it may ignore accident evolution and delay the safety systems, and neither case is exactly desirable. To mitigate the effects of thermal conduction and flashing phenomenon, this invention has each one of the pressure holes and its nearby space encased within a cylindrical coolant cladding to keep the temperature of the pressure-signal transmitting medium at safe condition. Experiments done in this laboratory prove that this local cooling strategy can effectively prevent loss of pressure-signal transmitting medium due to density gradient or vaporization. Thus, more than providing a method and a device to measure the void fraction of a gas/liquid two-phase stratified flow within a horizontal pipeline system, this invention provides a preventive method and a device to effectively overcome the problem of loss of measurement accuracy due to loss of pressure-signal transmitting medium encountered in transient flow system by the commonly used pressure-based instruments such as fluid-level meter, flowmeter, pressure-drop meter, pressure meter, etc."}
{"text": "Mesenchymal stem cells (MSCs) are capable of differentiating into a wide range of cell types including those associated with the bone development and maintenance (osteoblasts, osteocytes), cartilage, tendons, muscle, adipose, and more. The pluripotent nature of MSCs makes them a potentially valuable tool in therapeutics, particularly in the fields of disease (e.g. osteogenesis imperfecta), tissue repair and regeneration.\nThe molecular mechanisms underlying MSC differentiation are poorly understood. Microarray profiling studies have clearly identified genomic profiles that distinguish mother and daughter cell lineages, yet the triggers that orchestrate these processes have not been identified.\nRNA interference (RNAi) is a near-ubiquitous pathway involved in post-transcriptional gene modulation. A key effector molecule of RNAi is the microRNA (miRNA or miR). These small, non-coding RNAs are transcribed as primary miRNAs (pri-miRNA) and processed in the nucleus by Drosha (a Type III ribonuclease) to generate short hairpin structures referred to as pre-miRNAs. The resulting molecules are transported to the cytoplasm and processed by a second nuclease (Dicer) before being incorporated into the RNA Induced Silencing Complex (RISC). Interaction between the mature miRNA-RISC complex and messenger RNA (mRNA), particularly between the seed region of the miRNA guide strand (nucleotides 2-7) and regions of the 3′ UTR of the mRNA, leads to gene knockdown by transcript cleavage and/or translation attenuation."}
{"text": "The invention concerns a mass spectrometer for analysis of secondary ions and post-ionized neutral secondary particles with an ion source for creating a primary ion beam to irradiate a sample, and to produce secondary particles. The source possesses a heatable ion emitter that is coated in the area exposed to the field with a liquid-metal layer that contains an ionizable metal that is emitted and ionized as the primary ion beam. The primary ion beam contains metal ions with various stages of ionization and cluster statuses. In particular, the invention concerns a spectrometer unit for mass analysis of the secondary particles as well as the ion source of such a mass spectrometer.\nIn the description below, the conventional designation will be used for ions in clusters related to their mass and charge status, thus:Binp+wherein n=the quantity of atoms in a cluster, and p+=charge status.\nIt is known to use liquid metal sources in secondary-ion mass spectroscopy in particular when operated as time-of-flight secondary-ion mass spectroscopy (TOF-SIMS). Applicants have proposed a liquid metal gold-cluster ion source for a spectrometer (see prospectus: Liquid Metal Gold Cluster Ion Gun for Improved Molecular Spectroscopy and Imaging, published 2002) that represents the state of the art for the overall TOF-SIMS concept.\nThe efficiency of TOF-SIMS measurements with respect to primary ion beams from mono-atomic Gallium ions could be significantly increased using Gold Primary Clusters, e.g., of type Au3+. The disadvantage of the use of Gold as the material for the primary ion beam is that when Gold ions are created, those of type Au1+ predominate, while cluster formats such as Au2+ or Au3+ provide only low components of the overall ion current.\nBismuth has been used successfully during the intensive search for additional cluster-forming substances, containing only one natural isotope for secondary-ion mass spectroscopy. Bismuth is an an-isotopic element with a melting point of 271.3° C. Additionally, Bismuth alloys such as Bi+Pb, Bi+Sn, and Bi+Zn are known that possess lower melting points (46° C.-140° C.) than pure Bismuth. Pure Bismuth, however, is given preference for a liquid metal ion source.\nIn U.S. Pat. No. 6,002,128 it is noted that Bismuth is suited for the creation of charged particles. However, neither cluster formation nor the option of a liquid metal ion source with Bismuth is described. Also, the Japanese Patent No. 03-084435 proposes a calibration alloy for a secondary-ion mass spectrometer with which mass spectra with high resolution may be obtained. For this, the elements V, Ge, Cd, Os, and Bi are named as elements with high negative secondary ionization. The isotope patterns for the above-mentioned elements provide characteristic, repeatable spectra. However, this document does not mention cluster formation of a liquid metal ion source. Also, it is not indicated that Bismuth is well suited for cluster formation."}
{"text": "In current digital wireless systems, the traditional up-conversion chain (or significant portion thereof) is primarily analog and includes types such as super-heterodyne, low intermediate frequency (IF) and zero IF up-conversion technology. These technologies start with the conversion of inherently digital signals to analog signals through high performance digital-to-analog (D/A) converters, generally due to the higher frequencies involved. Once converted to the analog domain, various combinations of analog filters, amplifiers, mixers and modulators (and perhaps other analog elements) are cascaded to achieve the up-conversion from the output of the A/D converter(s) to the radio frequency (RF) band of interest (transmit RF signal).\nLikewise, on the receiver side, the traditional down-conversion chain (or significant portion thereof) is primarily analog including such types as super-heterodyne, low IF and zero IF down-conversion technology. To achieve the down-conversion, various combinations of analog filters, amplifiers, mixers and demodulators (and perhaps other analog elements) are utilized to achieve the conversion from the RF band of interest (receive RF signal) to the input to A/D converter(s).\nComponent variation, tolerances, and aging all affect the design requirements, costs, and manufacturability of the analog up-conversion (transmitter) and down-conversion (receiver) chains. Accordingly, there is needed a digital transmitter and digital receiver that utilizes digital technology for the up-conversion and down-conversion chains."}
{"text": "1. Field of the Invention\nThe present invention relates to an elliptical trainer, especially to an elliptical trainer that has two crank assemblies and allows a user to perform smooth movements.\n2. Description of the Prior Art(s)\nAn elliptical trainer is a stationary exercise machine that simulates walking or running without causing excessive pressure to the knee joints and decreases the risk of impact injuries. The elliptical trainer offers a non-impact cardiovascular workout, capable of varying from low to high intensity based on a resistance preference set by users.\nA conventional elliptical trainer has a base, a flywheel, a handle assembly, two linkage assemblies and a resistance assembly. The flywheel and the handle assembly are separately mounted'rotatably on the base. The linkage assemblies are respectively mounted on two'opposite sides of the base and connect pivotally to the handle assembly and eccentrically to the flywheel. Each linkage assembly has multiple links pivotally connecting to each other in series. The resistance assembly applies resistance to the flywheel. When the conventional elliptical trainer is operated, the linkage assemblies rotate the flywheel. However, as the links of the linkage assemblies swing or rotate, one or more toggle positions exist so that the linkage assemblies and the flywheel do not operate smoothly. Consequently, the conventional elliptical trainer is not able to smoothly simulate the walking or running movements.\nTo overcome the shortcomings, the present invention provides an elliptical trainer having two crank assemblies to mitigate or obviate the aforementioned problems."}
{"text": "This invention seeks to satisfy a need for a device which can monitor and indicate the state of the charge of a lead acid battery with increased accuracy and sensitivity. The monitor and indicator according to this invention is of the type which responds constantly and promptly to extremely minute changes in specific gravity of the battery liquid electrolyte during charging and discharging of the battery.\nHeretofore, specific gravity monitors and indicators for batteries have been hydrostatic systems which tend to be slow in their response to changes in specific gravity of the electrolyte. Furthermore, fluid movement in these hydrostatic indicating systems tends to be extremely slow because of reliance on slow diffusion for operation.\nAnother serious drawback in the prior art is the lack of any provision for excluding bubbles constantly generated in the electrolyte during charging from entering and clogging the indicator liquid column, thereby retarding, and, in some cases, completely disabling a movable indicator element in the column, such as a float.\nIn substantially completely eliminating the above and other known deficiencies of the prior art, the present invention provides a battery condition indicator which can be readily built into a commercial lead acid storage battery to continuously monitor and indicate the state of the charge on the battery in response to changes in specific gravity of the battery electrolyte over an extremely narrow range.\nThe improved accuracy and sensitivity of the indicating system for batteries is achieved by the utilization of a simple continuous flow arrangement for the liquid electrolyte, which includes an effective barrier to the induction of bubbles contained in the electrolyte into the electrolyte column or passage of the indicator, whereby such bubbles are prevented from clogging the indicator in whole or in part. In essence, the present invention makes use of a hydrodynamic indicator system as distinguished from the much more sluggish hydrostatic systems in the prior art.\nIn a preferred form, the invention utilizes an indicator vertical column for liquid electrolyte having at least a bottom bubble-excluding entrance. At least one buoyant or sunk float is contained in the electrolyte column between submerged spaced stops which are positioned relative to electro-optical emitterdetectors to enable a digital indicating response to the movement of the float to or from the buoyant and sunk positions following changes in specific gravity of the electrolyte over a very narrow range. The interruption of the optical path between emitter and detector elements of the indicator caused by the rise or fall of the float can produce a digital response through a state-of-the-art binary counting circuit connected with the emitter and detector elements and having a suitable display.\nVariations in the construction of the indicator and additional features and advantages thereof over the prior art will become apparent during the course of the following detailed description."}
{"text": "The present invention relates to detonators, and more particularly, to non-energetics-based detonators, detonator systems using non-energetics based detonators and methods of detonating explosives.\nIn various industries, such as mining, construction and other earth moving operations, it is common practice to utilize detonators to initiate explosives loaded into drilled blast holes for the purpose of breaking rock. For instance, commercial electric and electronic detonators are conventionally implemented as hot wire igniters that include a fuse head as the initiating mechanism to initiate a corresponding explosive. Such hot wire igniters operate by delivering a low voltage electrical pulse to the fuse head, causing the fuse head to heat up. Heat from the fuse head, generated in response to the low voltage electrical pulse, initiates a primary explosive, e.g., lead azide, which, in turn, initiates a secondary explosive output pellet, such as pentaerythritol tetranitrate (PETN) at an output end of the detonator. However, conventional hot wire igniters must rely on an extremely sensitive primary explosive to transition the detonation process from the fuse head to the corresponding explosive output pellet. Moreover, it is possible that the voltage and power requirements to function this type of conventional hot wire igniter may be encountered from inadvertent sources like static, stray currents and radio frequency (RF) energy.\nAnother exemplary detonator type is referred to as an exploding bridgewire detonator (EBW). The EBW includes a short length of small diameter wire that functions as a bridge. In use, explosive material beginning at a contact interface with the bridge wire transitions from a low density secondary explosive pellet to a high density secondary explosive pellet at the output end of the detonator. To initiate a detonation event, a high voltage pulse is applied in an extremely short duration across the bridge wire causing the small diameter wire to explode. The shockwave created from the bridge wire's fast vaporization initiates the low density secondary explosive pellet, such as PETN, which in turn initiates the high density secondary explosive pellet such as cyclotrimethylene trinitramine (RDX), at the output end of the EBW.\nYet another exemplary detonator type is referred to as an exploding foil initiator (EFI). A conventional EFI includes a thin metal foil having a defined narrow section. A polymer film layer is provided over the metal foil. To initiate a detonation event, a high voltage, very short pulse of energy is applied across the metal foil to cause the narrow section of the metal foil to vaporize. As the narrow section of the metal foil vaporizes, plasma is formed as the vaporized metal cannot expand beyond the polymer film layer. The pressure created as a result of this vaporization action builds until the polymer film layer is compromised, thus triggering a shock wave that initiates the detonation a connected explosive device."}
{"text": "1. Field of the Invention\nThe present invention generally relates to expandable endoprosthesis devices, in particular expandable intraluminal vascular grafts, generally called stents, which are adapted to be implanted into a patient's lumen, such as a blood vessel, to maintain the patency of the vessel. These devices are frequently used in the treatment of atherosclerotic stenosis in blood vessels, especially after percutaneous transluminal coronary angioplasty (PTCA) procedures, to prevent restenosis of a blood vessel. The present invention also relates to an expandable intraluminal vascular graft that can be used in any body lumen, and be used for drug delivery.\n2. Description of Related Art\nIn expandable stents that are delivered with expandable catheters, such as balloon catheters, the stents are positioned over the balloon portion of the catheter and expanded from a reduced diameter to an enlarged diameter greater than or equal to the artery wall, by inflating the balloon. Stents of this type can be expanded and held in an enlarged diameter by deformation of the stent (e.g., Palmaz U.S. Pat. No. 4,733,665), by engagement of the stent walls with respect to one another (e.g., Kreamer U.S. Pat. No. 4,740,207; Beck et al. U.S. Pat. No. 4,877,030; and Derbyshire U.S. Pat. No. 5,007,926), and by one-way engagement of the stent walls together with endothelial growth into the stent (e.g., Stack et al. U.S. Pat. No. 5,059,211)."}
{"text": "EP 0853515 B (CROWN CORK & SEAL TECHNOLOGY CORP.) is a blow forming apparatus in which a whole can body is held in a three part movable mould: base, centre and top part while compressed air is introduced to force the can body to conform to the interior profile of the moving mould parts. Whilst this is suitable for forming complex shapes, such shaping requires complex blow forming machinery and is not always be needed for holding a can when carrying out more simple rotational operations.\nU.S. Pat. No. 5,694,822 (REYNOLDS METALS COMPANY) describes an apparatus for trimming cylindrical can bodies. The can bodies are successively mounted with vacuum onto a support mandrel so that the open end portion of the can body projects between an outside circular cutting knife and an inside trimming edge. The invention of this patent is said to be less complex than traditional trimmers in which the closed end of the can body is held by vacuum onto a base pad. The base pad of the traditional trimmer comprises a holder which has a flat face and a circular centralising feature. The holder has holes or channels through which negative air pressure or “suction” is applied. This generally allows the can body to be held concentrically and permits rotational operations to be performed accurately in a plane perpendicular to the longitudinal axis of the article.\nAlthough considered by Reynolds to be over-complex, this method of holding can bodies or shallower cup-shaped articles is well suited to high-speed manufacturing equipment. In such equipment, containers are loaded onto and unloaded from holders at high speed. The application of suction is sometimes complemented by the application of small amounts of compressed air when releasing the container.\nThe base of a container such as a beverage can often has a bead or rim portion known as a stand (bead), which stiffens the container and allows it to rest on a flat surface without rocking. The container is thus be conveniently held rigidly by a holder using suction onto a flat surface and a seal is formed between the surface and the stand.\nIt is also known to provide a raised feature on the face of the holder, which fits within a recess in the base of the container and provides a circular centralising feature. If there is no suitable recess in the base of the container or cup-shaped article, the centralising feature comprises a raised feature, which fits outside the base of the article. However, machine motion often limits the ability to provide a holder of a larger diameter than the container, which constrains the design of the centralising feature.\nDue to the variability in manufacturing processes, it is often necessary to allow a radial clearance between the container or cup-shaped article and the centralising feature. Without this clearance, the container can fail to seat properly on the flat face, and this may impair the seal or cause it to rock on the holder. Unfortunately, this clearance can lead to non-concentrical holding of the container or article and allows sideways/radial movement of the article.\nIf seal or grip is lost due to improper seating of the container or cup-shaped article, then manufacturing defects can occur. For example, a cutting operation may not cut through the full circumference due to rotation of the container/article on the holder. Sometimes, the cut becomes spiral if the article is lifted away from the holder. If grip is lost in a cap or can flanging or curling operation, an incomplete or uneven flange or curl results. Similarly, if grip is lost in a lining operation, a defective lining is created.\nThis invention seeks to provide a holder, which addresses these problems. In the following text, which relates to the present patent application, expressions such as “container” or “cup-shaped article” are used interchangeably and are intended to cover a wide variety of items, including lids, closures and metal can bodies. Use of either expression is not intended to be limiting or to exclude equivalents or standard alternatives in the packaging trade."}
{"text": "The cellular components that comprise haematopoietic supportive marrow stroma and associated skeletal tissues are derived from a population of multipotential cells including multipotential mesenchymal stem cells (MSC) and their precursors (mesenchymal precursor cells (MPCs)). MPC reside within the bone marrow spaces located primarily in perivascular niches which surround blood vessels (Gronthos S et al., 2003 and Shi S et al., 2003). Over the last two decades, studies have focused on the regenerative potential of different human stromal/mesenchymal cell populations to reconstitute fat, bone, cartilage and muscle (Gronthos S et al., 2003 and Simmons P J et al., 1991). However, the successful application of these technologies is dependent on the ability to isolate purified populations of MPC and to subsequently manipulate their growth and differentiation ex vivo. Moreover, there is a need to overcome the various technical and safety concerns often associated with stem cell based therapies by assessing the safety and efficacy of MPC preparations in appropriate large pre-clinical animal models representative of human diseases. Ovine models of human orthopaedic and cardiac disease/trauma have been widely used, as sheep share similarities with human anatomy, physiology, immunology and embryonic development (Airey J A et al., 2004; Licchty K W et al., 2000 and Mackenzie T C et al., 2001).\nWhilst, several human multipotential cell specific markers such as STRO-1, CD106, CD146 have been described in the literature (Gronthos S et al., 2003 and Shi S et al., 2003), significant progress in examining the therapeutic potential of multipotential cells in ovine models of human disease has been limited due a lack of specific reagents which enable the isolation and characterisation of an equivalent multipotential cell population in ovine tissues. The development of biological reagents, such as monoclonal antibodies, reactive with both ovine and human multipotential cells would greatly facilitate the capacity to monitor and equate the functional properties of both ovine and human multipotential cell populations in pre-clinical and clinical trials, respectively.\nAny discussion of documents, acts, materials, devices, articles or the like which has been included in the present specification is solely for the purpose of providing a context for the present invention. It is not to be taken as an admission that any or all of these matters form part of the prior art base or were common general knowledge in the field relevant to the present invention as it existed before the priority date of each claim of this application."}
{"text": "Most conventional sensors for the detection of chemicals and/or fire are electrochemical sensors. Electrochemical sensors operate at a fixed potential. Changes in that fixed potential, caused by the presence of an analyte, signal the analytes presence. Most often these sensors are directed towards the recognition of a single analyte and utilize a single sensor or a simple array of sensors.\nCurrent concerns about homeland security, battlefield protection and terrorist activities have created an interest in sensors that detect a variety of chemicals, such as a variety of blood agents, toxic industrial chemicals, explosives, etc, preferably with specificity for the identification of type and quantity of chemicals detected.\nAlso, an interest has developed in fire prevention and fire detection. Along with fire detection, early detection of various gases or of precursors of fires, before the presence of flame and smoke, which are necessary for detection using conventional smoke detection/fire detection devices, is also desirable. For example, with increased automation on ships, the Navy has sought fire detection systems capable of improved performance over conventional smoke detectors. To the extent such precursors can be identified, fire prevention can occur in advance of mere fire detection to avoid damage to expensive infrastructure.\nAlthough the use of multiple conventional electrochemical sensors monitoring multiple criteria has been studied for this type of chemical identification with varying degree of success, most of the sensors did not have the specificity desired. For example, in fire detection, multiple criteria electrochemical sensors have difficulty discriminating fire-like nuisance sources, such as welding, grinding steel, and cutting with a torch. Further, multiple criteria electrochemical sensors are typically limited to only a few standardized test sources, without providing the range of detection capability often desired.\nCermet sensors are well known in the art that use various ceramic metallic (cermet) films. See U.S. Pat. No. 6,218,687, which is incorporated herein by reference in its entirety. Rather than operating based on a single potential, some cermet sensors characterize data using cyclic voltammetry, in which current is measured while manipulating the potential in a cyclical fashion. Thus, the data provided by such sensors is provided in a cyclical waveform.\nCermet sensor arrays have been used for more realistic and difficult conditions than prior electrochemical sensors and have been used to identify chemical analytes and fires, particularly distinguishing nuisance sources from harmful events. Such cermet sensors are capable of high temperature operation and can be fabricated using both thick and thin film techniques. Cermet sensors are small, lightweight and low cost alternatives to conventional electrochemical sensors. Plus, cermet sensors have design flexibility since different cermet sensors may respond to the same gas in different ways. Using conventional methods for data analysis of cyclic voltammetry, cermet sensors are capable of detecting very low level concentrations (part per billion) for a variety of analytes using aqueous electrolytes, but only higher levels (parts per million) for gas phase detection.\nHowever, the use of cyclic voltammetry creates complex data waveforms having a great deal of information, such as peak position, magnitude, and shape, all of which can be exploited for analytical purposes. Such waveforms contain data that incorporates variations in the presence, concentration, temperature, accessibility to the sensor, function of the sensor etc. which must be selectively analyzed before any meaningful information about the identity and/or quantity of a certain analyte or fire/non-fire condition can be ascertained.\nPreviously pattern recognition software has been used to analyze waveform data. However, the shear volume of data generated by cyclic voltammetry has caused concern about accuracy of probable identification of analyte sources, particular in the presence of quite similar nuisance sources. Further, the volume of data required may be more than a microprocessor embedded within a portable device used for field applications is capable of processing in a fast, real-time manner.\nWavelet analysis of analytical data prior to pattern recognition has been used in a variety of different applications. For example, wavelet analysis has been used in medicine to specifically look at identifying conditions from time-dependant data such as EEG and ECG scans and in identifying features of interest in medical imagery. Non-medical applications also included image analysis, as well as optical character recognition, and acoustic pattern recognition.\nIn sensing applications, wavelet analysis of linear voltammetric data of analytical chemical data, has been demonstrated as a means of de-noising linear voltammetric data. An “electric tongue” system has been developed for monitoring industrial liquid process streams. They utilize a “Wavelet Neural Network” as a pattern recognition technique which incorporates wavelet transformation into a feed forward multilayer neural network architecture. The parameters of the wavelet transform are thus modified by the error propagation structure of the neural network during training. Wavelet packet transforms have also been used for pattern recognition. The wavelet packet transforms have included the use of a Euclidean distance measure to generate the best wavelet basis set of an analytical signal for classification purposes. This approach has been generally applied only to spectroscopic and chromatographic data. Wavelet transformation has also been used for semi-conductor based gas sensors, particularly the use of discrete wavelet transform as a data pretreatment step prior to pattern recognition via a number of different techniques. Wavelet transform followed by classification by a PNN has been done by HPLC-DAD data in the classification of Chilean wines, peak detection in LC-TOF MS data and to capture features in NIR data.\nIn these various known data analysis techniques, the wavelet-transformed data are not down-selected or narrowed based on any particular classifying ability. As such, the data continues to include random variables and data points which are not necessary for classification purposes. The inclusion of unnecessary data points increases processing time and creates additional variables which can unnecessarily complicate or interfere with pattern recognition. In some cases, the wavelet coefficients may be selected based upon mere visual inspection. However, it is particularly difficult to discern narrow differences by mere visual inspection. In other cases, another “quick and dirty” method has been used to reduces the wavelet data, but not by any particular method particularly suitable for classification purposes. These “quick and dirty” methods do not ensure selection of the most desirable features of a wavelet, which can lead to lack of accuracy in pattern recognition."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of preventing retinal ganglion cell death, associated with glaucoma, by administering to retinal ganglion cells of a mammal, a compound which blocks the putative non-inactivating sodium ion channels of the above cell type.\n2. Brief Description of the Art\nGlaucoma is an optic neuropathy associated with elevated intraocular pressures which are too high for normal function of the eye, and results in irreversible loss of visual function. (See for example, Dreyer et al \"Elevated glutamate levels in the vitreous body of human and monkeys with glaucoma\", Arch. Ophthalmology 114:299-305, 1996) It is estimated in medical science that glaucoma afflicts approximately 2 per cent of the population over the age of forty years, and is therefore a serious health problem. Ocular hypertension, i.e. the condition of elevated intraocular pressure, which has not yet caused irreversible damage, is believed to represent the earliest phase of glaucoma. Many therapeutic agents have been devised and discovered in the prior art for the treatment or amelioration of glaucoma and of the condition of increased intraocular pressure which precedes glaucoma.\nPrimary open angle glaucoma (POAG) is associated with a rise in intraocular pressure (IOP). This increase in IOP is believed to contribute to the loss of optic nerve function which ultimately leads to blindness. Reduction of IOP is therefore a crucial component in the management of POAG. However, in many individuals lowering of IOP is not sufficient or ineffective in preventing vision loss associated with POAG.\nIt is thought that a novel class of sodium channels residing within the optic nerve of the rat are responsible for damage to the rat optic nerve following anoxia or hypoxia. However, in glaucoma the sequence of pathological events leading to the loss of optic nerve function, is not known.\nThe drugs currently utilized in the treatment of glaucoma include miotics (e.g., pilocarpine, carbachol, and acetylcholinesterase inhibitors), sympathomimetrics (e.g., epinephrine and dipivalylepinephrine), beta-blockers (e.g., betaxolol, levobunolol and timolol), alpha-2 agonists (e.g., para-amino clonidine) and carbonic anhydrase inhibitors (e.g., acetazolamide, methazolamide and ethoxzolamide). Miotics and sympathomimetics are believed to lower intraocular pressure by increasing the outflow of aqueous humor, while beta-blockers, alpha-2 agonists and carbonic anhydrase inhibitors are believed to lower intraocular pressure by decreasing the formation of aqueous humor. All five types of drugs have potential side effects. Miotics, such as pilocarpine, can cause blurring of vision and other visual side effects which may either decrease patient compliance or require termination of miotic drug therapy. Carbonic anhydrase inhibitors can also cause serious side effects which affect patient compliance and/or necessitate withdrawal of the drug therapy. At least one beta-blocker, timolol, has increasingly become associated with serious pulmonary side effects attributable to its effect on beta-2 receptors in pulmonary tissue.\nAs a result additional antiglaucoma drugs are being developed, e.g., prostaglandin derivatives, muscarinic antagonists, etc. However, none of the above drugs are designed to directly interact with the retinal ganglion cell and its associated axon.\nThus, it would be desirable to prevent the loss of ganglion cell body and axon function, which may be associated with glaucoma by a biological mechanism which does not modulate aqueous humor dynamics and therefore intraocular pressure. Moreover, it would be desirable to treat the retinal ganglion cell body and axon of a mammal directly to prevent the destruction thereof by the glaucomatous condition."}
{"text": "The management of rewards for an account is known, for example, as disclosed in U.S. Published Patent Applications 2002/0087410 (Walker et al.), incorporated by reference herein. Unfortunately, such programs are not self-learning.\nThus, there is a long-felt need to provide a system and a method to manage rewards for a multi-tiered account that is dynamic and can be readily adapted to meet various and variable requirements."}
{"text": "Vehicles are used both in conditions with good visibility, such as during the day, and in conditions with poor visibility, such as at night or in bad weather. Thus, in conditions where the visibility is good, the driver can see people and objects such as fallen objects around the vehicle by looking at the actual landscape with his or her (hereinafter, the driver will be referred to in the masculine form for simplicity) own eyes. In conditions where the visibility is poor, however, it is difficult for the driver to see all of these by looking the landscape. Hence, vehicle surroundings information output systems have been developed which inform the driver of the existence an object about which the driver is to be warned (hereinafter also referred to simply as “object”), such as a person or object around the vehicle.\nJapanese Patent Application Publication No. JP-A-2003-291688, for example, describes one kind of vehicle surroundings information output system which is a so-called head-up display (HUD) that displays the landscape in front of the vehicle on the front windshield as both an object about which the driver is to be warned and a virtual image.\nWhen an object is detected in front of the vehicle, the vehicle surroundings information output system (i.e., the head-up display system) described in Japanese Patent Application Publication No, JP-A-2003-291688 displays information (a solid star symbol) of strong visual stimulus to the driver and moves information (a triangular outline) to guide the driver's line of sight from that place to the object on the head-up display. That is, after an object is detected, the vehicle surroundings information output system described in Japanese Patent Application Publication No. JP-A-2003-291688 notifies the driver of the existence of that object by displaying information that stimulates the driver's vision (i.e., catches the driver's eye) and draws his line of sight to it (i.e., by displaying highly visually alluring information).\nAlso, other similar vehicle surroundings information output systems (head up display systems) which display such highly visually alluring information are described below.\nJapanese Patent Application Publication No. JP-A-7-61257 describes a vehicle surroundings information output system which displays highly visually alluring information when the driver is not looking at the display of the head-up display and displays not so visually alluring information when the driver is looking at the display from the head-up display.\nJapanese Patent Application Publication Nos. JP-A-2004-30212 and JP-A-2001-357498 describe vehicle surroundings information output systems which suppress the provision of information relating to a detected object as necessary. More specifically, Japanese Patent Application Publication No. JP-A-2004-30212 describes a vehicle surroundings information output system which prevents the driver from receiving more information than is necessary, thereby reducing the burden on the driver, by providing the minimum necessary information relating to objects while suppressing the provision of highly visible and visually alluring information relating to objects. Meanwhile, Japanese Patent Application Publication No. JP-A-2001-357498 describes a vehicle surroundings information output system which prevents the driver from receiving more information than is necessary, thereby reducing the burden on the driver, by suppressing information relating to objects near where the driver is looking.\nIncidentally, Japanese Patent Application Publication No. JP-A-2000-251200 discloses an obstacle detection system structured such that warnings are issued more often when visibility of detected objects is poor than when it is good.\nHowever, the vehicle surroundings information output system described in Japanese Patent Application Publication No. JP-A-2003-291688 provides the driver with high visually alluring information instantaneously after an object has been detected. That is, the head-up display displays highly visually alluring information when it is still unclear as to whether that object is real or not. Therefore, even if that object does not actually exist, the driver's eyes are still drawn to that highly visually alluring information. When the driver wants to drive by detecting objects on the actual landscape with his own eyes, his line of sight is forced from the actual landscape by the false information. The same can also be said for the vehicle surroundings information output system described in Japanese Patent Application Publication No. IP-A-7-61257 when highly visually alluring information is displayed when the driver is not looking at the head-up display.\nIn this case, even though these kinds of problems exist, if no information is provided to the driver until it is clear as to whether the object is real or not and that object happens to really exist, and further, if a warning has to be issued quickly, the notification to the driver of that object ends up being delayed. The same may also be said for the vehicle surroundings information output systems described in Japanese Patent Application Publication Nos. JP-A-2004-30212 and JP-A-2001-357498 when the provision of information regarding a detected object is suppressed as necessary because the driver feels it is annoying.\nThat is, despite the fact that some drivers want to drive giving priority to information obtained from the actual landscape while other drivers want to drive giving priority to information provided by the vehicle surroundings information output systems described above, these systems provide and suppress information relating to objects with disregard to the intentions of both kinds of drivers."}
{"text": "Various data communication schemes are available for radio communication systems. Modulation techniques (e.g., analog or digital modulation) may be utilized in such communication schemes. In addition, encoding and decoding processes may also be utilized to improve the signal integrity of the data being communicated.\nVarious data transmission schemes are available for information systems. Cipher techniques may be utilized in such transmission schemes. In addition, encoding and decoding processes may also be utilized to improve the signal integrity of the data being communicated."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor devices, and, more particularly, to semiconductor devices that combine high-speed performance and low power consumption.\n2. Description of the Related Art\nAn example of reducing power consumption by substrate bias control is described in 1996 IEEE International Solid-State Circuit, Digest of Technical Papers, (1996), pp. 166-167.\nWith the recent popularity of low-power CMOS LSIs (Complementary Metal Oxide Semiconductor Large Scale Integrated circuits), a trend has developed to maintain high-speed operation by decreasing the threshold voltage VT of the MOSFETs as the operating or supply voltage is dropped. When the supply voltage drops to 2 V or lower and when the threshold voltage VT is decreased to 0.5 V or lower correspondingly, however, the subthreshold leakage current increases, whereby the transistor cannot be cut off completely. Consequently, the standby current of the LSI chip increases, which represents a bottleneck in the design of a system that includes a battery-powered CMOS LSI chip. Furthermore, the current during normal operation also increases as the threshold voltage VT increases.\nIn order to break the bottleneck, a well-known system achieves a high-speed operation by decreasing the threshold voltage of each of the MOSFETs in the chip during normal operation, and decreases the standby current by increasing the threshold voltage at the time of standby. Nevertheless, the following three problems exist in this system:\n(1) An overcurrent flows because of latching-up when the power supply is turned on, and the wiring in the CMOS LSI chip may fuse, or the normal supply voltage may become inapplicable as the load exceeds the current capacitance of the power supply. This problem is caused because the layout and connections of the circuit are designed so that the substrate (well) and source of the MOSFET are not at equipotential.\nFor example, when a p-channel MOSFET (PMOSFET) is used for applying a positive supply voltage (e.g., 1.8 V) to the source (p-layer), the pn junction between the source and well is excessively biased in the forward direction because the well (n-well) remains at a floating 0 V just until the application of the supply voltage, thus causing latch-up of the CMOS. In the case of conventional CMOS LSI products at 2 V or higher, the pn junction is never biased in the forward direction as in the normal operation thereafter, even during the application of the supply voltage, since the well and source of the MOSFET are connected so that both are at equipotential as much as possible. Since the threshold voltage VT is constant at all times at a value of substantially 0.5 or higher, moreover, there is no problem of subthreshold current.\nIn the case of an n-channel MOSFET (NMOSFET), the problem is not so serious. When the supply voltage is applied to the drain, the substrate (p-well) of the NMOSFET is at a floating 0 V and the source is fixed to an earth potential of 0 V, because the pn junction, formed between the drain and the well, is not biased in the forward direction. However, there is a subthreshold current flowing between the drain and source when the threshold voltage is 0.5 V or lower. By separately controlling the well and source, the threshold voltage in the CMOS LSI is lowered.\n(2) The time required to switch the normal mode to the standby mode and the time required to switch the standby mode to the normal mode are extremely long, on the order of xcexcs. Assuming that the substrate voltage is generated on-chip, by a charge pumping circuit for pumping the capacitor in the chip, the output current is limited to a low level. On the other hand, the transistor in the chip is used to connect the power supply terminals of the substrate in common, and consequently the total substrate capacitance has an extremely large value (100 pF or greater). Therefore, a large load (substrate) capacitance is driven by a substrate-voltage generating circuit whose current driving capability is low when the mode is switched, so that the response time tends to become longer.\n(3) The subthreshold current flows everywhere, even in the CMOS circuit, thus increasing the operating current of the whole chip. This problem exists because, in the inactive state, the threshold voltage of the transistor in the CMOS circuit or circuit block is low during normal operation.\nThe present invention controls the substrate or well voltage of a transistor in a manner that solves the foregoing three problems.\nAn object of the invention is to hinder the latching-up that occurs when the supply voltage is applied to a CMOS circuit that includes MOSFETs having a low threshold voltage, or when the supply voltage is cut off to the CMOS circuit.\nAnother object of the invention is to decrease the subthreshold current during normal operation.\nA further object of the invention is to realize low power consumption while maintaining high operating speed, for a CMOS circuit that operates at a voltage of 2 V or below, a CMOS LSI, and a semiconductor device using the CMOS circuit.\nIn general, the invention achieves these and other objects by controlling the well voltage of a CMOS circuit when the power supply is turned on and when cut off, and during operation.\nIn one embodiment of the invention, in which the CMOS circuit includes MOSFETs that cannot be cut off substantially satisfactorily during normal operation, after the well voltage is applied to the well of the CMOS circuit so that the MOSFETs can be cut off, the supply voltage is applied to the CMOS circuit.\nIn another emboidment, after a third supply voltage (generated from a first supply voltage by a voltage conversion circuit) is applied as a well voltage to the well of the CMOS circuit, a second supply voltage is applied to the CMOS circuit.\nIn yet another embodiment, the invention provides a circuit for fixing the well potential of the CMOS circuit, and a circuit for varying the well potential of the MOSFETs by capacitive coupling according to the variation of the input signal of the CMOS circuit.\nIn still another embodiment, the invention provides a semiconductor device that includes a dynamic memory cell comprising a MOSFET, a capacitor, and a CMOS circuit, wherein the well potential of the MOSFETs constituting the CMOS circuit is subjected to a pulse variation, and wherein the substrate voltage of the dynamic memory cell is substantially a DC supply voltage.\nIn another embodiment, the invention provides a semiconductor device including a static memory cell that is operated at a high voltage and is constituted by MOSFETs of high threshold voltage, and a CMOS circuit operated at a low voltage and constituted by MOSFETs having a low threshold voltage. The well potential of the MOSFETs constituting the CMOS circuit is subjected to a pulse variation.\nIn yet another embodiment, the invention provides a semiconductor device including at least one CMOS circuit, a standby control circuit, and a voltage conversion circuit, wherein the voltage generated by the voltage conversion circuit is supplied to the standby control circuit, and the standby control circuit varies the well potential of the CMOS circuit using the output of the voltage conversion circuit, depending on operating conditions. A capacitor having a capacitance that is greater than the capacitance of the well is connected to the output of the voltage conversion circuit."}
{"text": "This invention relates to computer systems and circuits; and more particularly, to methods and apparatus for generating an appropriate reference voltage level for differential signaling.\nComputer systems are consistently being updated to operate at faster rates. In doing so, many of the components of these system must operate within new parameters, or be replaced at an additional cost. Achieving these higher operating rates further requires devices within these systems to communicate among themselves at faster rates.\nA technique employed to increase the rate at which two devices can communicate with each other is decreasing the voltage swing between a communication signal\"\"s high and low values. Typically, this includes reducing the supply voltage level at which these devices operate. This decrease in voltage swing allows a signal to transition faster between values. However, as this voltage swing is decreased, additional precautions must be taken to ensure that noise does not interfere with the signal.\nDifferential signaling is one method used in communicating between devices in a system and is useful in avoiding errors induced by noise. Typically, each sending device provides a voltage reference and a communication signal which are typically connected to a differential amplifier. Because approximately the same noise is typically induced on both the voltage reference and a communication signal, a receiving differential amplifier produces a relatively xe2x80x9ccleanxe2x80x9d communications signal.\nHowever, not all sending devices can provide an appropriate reference voltage. For example, when a system is upgraded or a new version is introduced that requires a different reference voltage level and possibly uses a different supply voltage, an older component of the system might not be able to provide the new reference voltage.\nSuch is the problem introduced by the Accelerated Graphics Port (xe2x80x9cAGPxe2x80x9d) Interface Specification Revision 2.0, May 1998, which requires a universal AGP target (such as the 82465GX GXB manufactured by Intel Corporation) to support both 3.3 Volt and 1.5 Volt operation on the AGP interface supply voltage (Vddq). The specification stipulates a voltage reference that is nominally 0.5*Vddq for 1.5 Volt operation, and nominally 0.4*Vddq for 3.3 Volt operation. Therefore, the value of the reference voltage must be adjusted to match the requirements based on the supply voltage used. What is needed is a system for efficiently providing the appropriate reference voltage.\nAccording to the invention, systems, apparatus and methods are disclosed for providing one or more bimodal reference voltages for using in differential signaling. Typically, a system comprises a first component requiring a first reference voltage; a second component requiring a second reference voltage; a power source to produce a supply voltage; a first voltage reducer electrically coupled to the power source, the first voltage reducer to produce a third reference voltage using the supply voltage; a multiplexor electrically coupled to the first voltage reducer to produce the first reference voltage by selecting between at least the third reference voltage and a fourth reference voltage; and a second voltage reducer electrically coupled to the power source and to the second component, the second voltage reducer to produce the second reference voltage using the supply voltage."}
{"text": "As communication technology develops, cell phones and other mobile terminals have become more powerful. Not only can one make telephone calls and send short messages with a mobile terminal, one can also install various applications (apps) on the cell phones and the other mobile terminals to perform various functions. When using a mobile terminal, the mobile terminal often captures and stores an image within an application (app). The capturing and storing of the image process is also known as screen capture.\nThe screen capture technique of existing mobile terminals using the Android system includes executing the screen capture through SurfaceView. However, images captured through the SurfaceView screen capture will turn black within the SurfaceView zone due to the way screen capture in SurfaceView is implemented. Because SurfaceView and commonView are in different hierarchies with respect to rending, only black images are captured through the SurfaceView screen capture technique when using a mobile terminal. The screen capture result is poor, and only a single method for implementing the screen capture exists."}
{"text": "This application is based upon and claims the benefit of priority from the prior Japanese Patent Application No. 2000-246913, filed Aug. 16, 2000, the entire contents of which are incorporated herein by reference.\nThis invention relates to a method of manufacturing a synthesis gas to be employed for the synthesis of gasoline, methanol or dimethyl ether by way of the GTL (Gas to Liquid) process.\nA synthesis gas comprising hydrogen (H2) and carbon monoxide (CO) has been employed as a raw material for the synthesis of gasoline, etc. by way of the GTL (Gas to Liquid) process according to the Fisher-Tropsch reaction system.\nThis synthesis gas has been conventionally manufactured by a method wherein steam and carbon dioxide are added at first to a natural gas employed as a raw gas to prepare a mixed gas, which is then fed to the reaction tube of reformer which has been heated to a predetermined temperature so as to steam-reform the natural gas together with carbon dioxide, thereby manufacturing the synthesis gas containing hydrogen (H2) and carbon monoxide (CO).\nBy the way, depending on the place of origin, the natural gas occasionally contains hydrogen sulfide (H2S) and carbon dioxide (for example, CO2: 7.1% by volume, and H2S: 0.6% by volume). When such a natural gas is to be employed as a raw material, it is usually practiced to remove hydrogen sulfide prior to the step of adding steam and carbon dioxide to the natural gas. In this case, the removal of hydrogen sulfide from the natural gas is conventionally performed by means of amine absorption method.\nHowever, if it is desired to remove hydrogen sulfide to a level of the order of ppm by means of the amine absorption method, carbon dioxide included in the original natural gas is also removed together with hydrogen sulfide. As a result, it becomes impossible to utilize the carbon dioxide that has been originally included in the natural gas, so that a large quantity of carbon dioxide is required to be supplied to the natural gas prior to the stage where the natural gas is to be transferred to the reaction tube of the reformer, thus increasing the manufacturing cost of the synthesis gas.\nTherefore, an object of this invention is to provide a method for manufacturing a synthesis gas, which makes it possible to effectively utilize the carbon dioxide in a natural gas containing hydrogen sulfide and carbon dioxide by selectively removing only the hydrogen sulfide therefrom in the synthesis gas manufacturing method employing a reformer, thereby making it possible to reduce the quantity of carbon dioxide to be added to the natural gas prior to the stage where the natural gas is to be transferred to the reformer.\nNamely, this invention provides a method of manufacturing a synthesis gas containing hydrogen and carbon monoxide, which comprises steps of;\nremoving only hydrogen sulfide from a natural gas containing hydrogen sulfide and carbon dioxide by permitting the natural gas to pass through a hydrogen sulfide-removing device filled with a hydrogen sulfide absorbent;\nadding carbon dioxide and steam to the natural gas which the hydrogen sulfide has been removed to prepare a mixed gas; and\nfeeding the mixed gas into a reaction tube of a reformer, thereby permitting mainly a steam reforming reaction to take place in the mixed gas.\nThe method of manufacturing a synthesis gas according to this invention may be performed in such a way that before the natural gas is fed to the hydrogen sulfide-removing device, the natural gas is forced to pass through a convection portion communicated with a combustion radiation portion of the reformer, thereby heating the natural gas up to a temperature which is suited for the reaction between the hydrogen sulfide in the natural gas and the hydrogen sulfide adsorbent.\nThe method of manufacturing a synthesis gas according to this invention may be performed in such a way that the carbon dioxide recovered from a combustion exhaust gas generated at the combustion radiation portion of the reformer is utilized as a carbon dioxide source.\nThe method of manufacturing a synthesis gas according to this invention may be performed in such a way that the carbon dioxide recovered from the synthesis gas at a downstream side of the reformer is utilized as a carbon dioxide source.\nIt is preferable in the method of manufacturing a synthesis gas according to this invention that the hydrogen sulfide adsorbent is at least one oxide selected from triiron tetraoxide (Fe3O4) and zinc oxide (ZnO).\nIt is preferable in the method of manufacturing a synthesis gas according to this invention that the hydrogen sulfide-removing device is provided with at least one unit of first desulfurizing column filled with the hydrogen sulfide adsorbent comprising triiron tetraoxide and with a second desulfurizing column filled with the hydrogen sulfide adsorbent comprising zinc oxide, and that the natural gas containing hydrogen sulfide and carbon dioxide is permitted to successively pass through one column selected from these first desulfurizing columns and second desulfurizing column.\nIt is preferable that the hydrogen sulfide-removing device is provided with at least three units of first desulfurizing columns to be filled with triiron tetraoxide, wherein a first placed first desulfurizing column among the first desulfurizing columns is designed to execute an adsorption operation of hydrogen sulfide, a second placed first desulfurizing column among the first desulfurizing columns is designed to execute an operation of regenerating the adsorbent (iron sulfide) on which hydrogen sulfide is adsorbed, and a third placed first desulfurizing column among the first desulfurizing columns is designed to execute an operation of reducing the adsorbent that has been regenerated, these operations being sequentially executed.\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out hereinafter."}
{"text": "F125 is an HFC (hydrofluorocarbon) which can be employed as substitute for F115 (a CFC or chlorofluorocarbon) in the field of low-temperature refrigeration.\nKnown processes for obtaining F125 include, inter alia, the fluorination of perchloroethylene or its fluorinated derivatives like 1,1-difluoro-l,2,2-trichloroethane (F122), 1,1-dichloro-2,2,2-trifluoroethane (F123) and 1-chloro-1,2,2,2-tetrafluoroethane (F124), the fluorination of chlorotrifluorethylene (F1113) or the chemical reduction of F115, especially the hydrogenolysis of the latter. In most cases these synthesis routes to F125 produce a crude F125 which is contaminated by significant quantities (several hundred ppm to several tens per cent) of F115, either through formation of byproducts due to the high temperatures needed to obtain high yields of F125 (in the case of fluorination reactions), or because of nonquantitative conversions of the starting material (in the case of the hydrogenolysis of F115).\nAs indicated in U.S. Pat. No. 5,087,329, the separation of the compounds F125 and F115 by distillation is very difficult, or even impossible if it is desired to obtain a very low F115 content.\nNevertheless, there is at present a trend towards a commercial F125 of high purity in which the F115 content might be of the order of a hundred ppm or so. Because of the difficulties encountered in separating the compounds F125 and F115, obtaining such low F115 contents does not appear to be feasible with a distillation line of realistic size. Processes for purifying F125 which make it possible to gain access to these very low F115 contents are therefore being investigated at present.\nThe process for separating F125 and F115 which is described in the abovementioned U.S. Pat. No. 5,087,329 relies on extractive distillation consisting in adding to the F125-F115 mixture before distillation an extractant consisting of an optionally hydrogenated and/or chlorinated C.sub.1 -C.sub.4 fluorinated hydrocarbon which has a boiling point of between -39.degree. C. and +50.degree. C. As examples of such extractants the patent mentions 1,2-dichlorotetrafluoroethane (F114), 1,1-dichlorotetrafluoroethane (F114a), 1,1,2-trichlorotrifluoroethane (F113), 1,1,1-trichlorotrifluoroethane (F113a), 2-chloro-1,1,1,2-tetrafluoroethane (F124), 2,2-dichloro-1,1,1-trifluoroethane (F123), trichlorofluoromethane (F11) and octafluorocyclobutane (F-C318).\nIn Patent Application EP 508,631, which describes the production of HFC compounds by liquid-phase chemical reduction of chlorinated, brominated or iodinated compounds with a metal hydride or a complex of such a hydride, it is indicated that this process may be of interest for purifying certain HFCs like F125. The purification of F125 is not exemplified. On the other hand, example 6 describes the reduction of F115 with LiAlH.sub.4 in diglyme; the best result indicates a 60% conversion of F115 to F125."}
{"text": "The present invention relates to a process for waterproofing leather, particularly leather for the manufacture of shoes, clothes or leather accessories. The present invention also relates to the leather obtained by said process.\nLeathers, particularly those which are used for the manufacture of shoes, clothes or leather accessories, are known to be waterproofed by chemical treatments consisting in applying thin layers of water-repellent substances outside the leather itself, for example by means of sprayers. Because of the low efficiency of this process in the long term, a new process has been conceived, consisting in sewing inside the leather a fabric lining combined with a semi-permeable film, which not only prevents water from further penetrating into the shoe or dress article, but also allows the body outward transpiration.\nHowever, this process is disadvantageous in that it allows water to penetrate anyway under the leather, and therefore, especially if the latter is used for the shoe-manufacture, an undesired water-cushion is formed between the waterproofing fabric and the leather internal surface. Further, this process necessarily requires the application of a lining combined with a semi-permeable film, which, besides increasing the production costs, is in some cases inadvisable, for example in the manufacture of summer shoes and clothes. U.S. Pat. No. 5,598,644 discloses a process for waterproofing a tread of leather, wherein a semi-permeable membrane is glued to the perimetric regions thereof. However, this process does not avoid the formation of water-cushions and the need of an upper perimetric part for fastening the membrane to the tread, thus reducing its transpiration properties."}
{"text": "Integrated circuits (“ICs”) such as, for example, programmable integrated circuits can be programmed to perform specific functions. One example of programmable integrated circuits is programmable logic devices (“PLDs”) which can be programmed by a user to perform specified logic functions. One type of programmable logic device, called a field programmable gate array (“FPGA”), is popular because of a superior combination of capacity, flexibility, time-to-market, and cost.\nIn a typical architecture, the FPGA includes a processor, an array of configurable logic blocks (“CLBs”), and programmable input/output blocks (“IOBs”) which provide external access to the CLBs. The processor, CLBs, and IOBs are interconnected by programmable fabric routing. The fabric routing comprises many interconnect wires and associated programmable interconnect points (PIPs) which can be programmed to selectively route signals among the various CLBs and IOBs.\nThe processor typically implements numerous debug features such as a trace feature. When the trace feature is enabled, the instruction stream executed by the processor is monitored and certain debug events are used to cause trigger events. When a trigger event occurs, the processor asserts an output-trigger-event signal to indicate that the trigger event occurred. The output-trigger-event signal is typically modified outside the processor and fed back to the processor using the fabric routing.\nFor proper operation of the trace feature, there is typically a constraint that the output-trigger-event signal be fed back to the processor within less than a fraction of one clock period. However, placement of pins on the processor is based on optimizing processor performance rather than optimizing debug features, and thus the debug pins associated with the input and output trigger-event signals are typically placed far apart from each other. Thus, the delays associated with using the fabric routing (e.g., the delay associated with driving the output-trigger-event signal to two or more CLBs) are too great to satisfy this constraint when the processor is running at or near maximum speed. As a result, features such as trace collection cannot be done reliably unless the processor is configured to run at a slower speed. However, processor performance and efficiency are significantly decreased if the processor is configured to run at a slower speed. In addition, a problem may not appear if the processor is configured to run at a slower speed because the problem may only occur when the processor is running at the faster speed. For example, deadlock may occur between the processor and a hard disk drive only when the processor is running at full speed. However, deadlock may not occur when the processor is running at a slower speed.\nTherefore, it is desirable to feed back a signal such that stringent timing requirements are satisfied while still maximizing processor performance."}
{"text": "The present invention relates to a method of refreshing a dynamic RAM (Random Access Memory) used by any personal computer, more particularly to a new method of refreshing a dynamic RAM so that the desired refreshing of a dynamic RAM can be executed by means of relatively simple circuitry, and therefore any of the high-speed devices using the real-time means such as a floppy disc driver unit can easily be accessed.\nThe semi-conductor manufacturing technology has achieved an overall development in recent years. This modern technology has successfully brought a variety of ICs into the field of the computer memory means. The problem of trading cost economy for further improvement of memory density potentials has been solved by now, and as a result, a large quantity of low-cost ICs, having greater capacities, have been delivered to the industrial markets, including ROM, static RAM, dynamic RAM, and others. Of these, compared to the static RAM, the dynamic RAM must be periodically subject to a refreshing process in order to properly store a variety of data in memory. Usually, such a periodic refreshing process must be executed at intervals of every 2 microseconds. In terms of the operational speed, such a dynamic RAM is rather disadvantageous when compared to other memories having identical access time such as the static RAM. This is because a short period of delay must be taken into account covering a period from the entry of the address up to the data output. Conversely, the dynamic RAM can be advantageously used in a system such as one incorporating a high-density memory, for example, that incorporates a minimum of 16K byte memory. Any dynamic RAM stores its bit data in the bit-cell capacitor as charge. A typical construction of the dynamic RAM is shown in FIG. 1 of the attached drawings. During the write mode, any desired data is fed to the input of the dynamic RAM so that switch S1 is closed. Capacitor C is either charged or discharged according to the status of the input data. During the read mode, switch S2 is closed so that the voltage in capacitor C is compared to the reference voltage Vref of comparator CP, and then according to the difference of the voltages compared, either \"1\" or \"0\" of the binary code is transmitted to the output of the dynamic RAM. When the dynamic RAM remains in the data-hold mode, all the switches S1, S2, and S3 remain open so that the data can remain in capacitor C. Nevertheless, since the MOS capacitor is characteristically subject to leakage resistance R, the capacitor charge will be gradually discharged. To compensate for this, a process called \"refreshing\" must be applied to any dynamic RAM being used. To achieve this, the dynamic RAM controller must be subjected to the \"refreshing\" process at intervals of every 2 microseconds so that the data stored in the bit-cell can properly be held unaffected.\nConventionally, in regard to the refreshing process applied to the dynamic RAM, a variety of means are made available, for example, \"burst mode refreshing process\" capable of refreshing all lines within a specific period of time, for example: within 2 microseconds, and the other one called \"cycle steal refreshing process\" (or called \"single cycle refreshing process\") capable of sequentially refreshing every line at specific intervals. Of these, the \"burst mode refreshing process\" can be executed by means of a relatively simple circuitry. However, since this process requires entire operation of the CPU (central processing unit) to periodically halt for a relatively long period, for example, at least for a period of 50-60 microseconds for each cycle, and as a result, using this process, any real-time high-speed devices such as a floppy disc driver system cannot be accessed. Conversely, since the \"cycle steal refreshing process\" uses the surplus time of the CPU bus cycle, this process is advantageous because it does not cause the CPU to actually decrease the speed needed to execute any operation. On the other hand, this second process still has a difficult problem in that a considerably high speed is needed to properly operate the circuits, and the operation timing cannot easily be controlled. Of these two refreshing processes, in reference to the attached drawings, some of the existing problems related to the \"burst mode refreshing process\" are described below.\nA simplified block diagram shown in FIG. 2 and its flowchart shown in FIG. 3 respectively denote a significant disadvantage of the conventional \"burst mode refreshing process\". The data contents of a dynamic RAM (DRAM) loaded in the CPU, of a microprocessor are continuously transmitted to the floppy disc controller (FDC). If the time at which the data contents were loaded from the DRAM to the CPU exactly coincides with the period when the DRAM is being refreshed by the \"burst mode refreshing process\", the CPU will be obliged to remain in a stand-by mode for a long time, and as a result, the contents of the dynamic RAM loaded into the CPU cannot be written into the floppy disc controller (FDC) within a pre determined period. In FIGS. 2 and 3, symbol DRAM represents a dynamic RAM, symbol FDC represents a floppy disc controller, symbol FDD represents a floppy disc device, and symbol CPU represents the central processing unit. Symbol DRQ is a data request signal output from the FDC to the CPU, while symbols RD and WT respectively represent the read and write signals output from the CPU to the FDC.\nIn reference to the flowchart of the CPU operation shown in FIG. 3, operations of the circuit shown in the block diagram are described below. During step S1, if the floppy disc controller (FDC) outputs a data request signal (DRQ=High) to the CPU, the mode then proceeds to step S2 where the contents of the address designated dynamic RAM (DRAM) are loaded into the CPU. After being temporarily stored in the CPU, said contents are then sent to the floppy disc controller (FDC) during the next step S3. This operation is continuously executed until all the contents of the dynamic RAM (DRAM) loaded in the CPU are completely transmitted to the floppy disc controller (FDC) during the next step S4. In this case, such a floppy disc controller can consist of a microprocessor. If any of the single side floppy discs is used as the floppy disc controller (FDC), all the data contained in the dynamic RAM (DRAM) must be fed to the floppy disc controller within 32 microseconds from the time the data request signal (DRQ) is received. This corresponds to the processing period covering steps S1 through S3 shown in the block diagram of FIG. 2. These processes are actually executed within about 30 microseconds. Using a programmed system like this, when transmitting the data contents of the dynamic RAM (DRAM) temporarily stored in the CPU to the floppy disc controller (FDC), a when such contents output from the dynamic RAM are being fed to the CPU, if this timing coincides with the period of the \"burst mode refreshing process\", the CPU will be obliged to stand by for a long time, and as a result, the necessary data stored in the CPU cannot be written into the floppy disc controller within the predetermined period of time. To prevent this, if the refreshing process is stopped, the period between refreshes of the RAM which should normally last 2 microseconds will run out, and as a result, the data contents remaining in the dynamic RAM will eventually be destroyed."}
{"text": "a) Field of the Invention\nThe invention is directed to a front sole holding device according to the preamble of claim 1."}
{"text": "1. Field of the Invention\nThe present invention relates to a solar battery controller.\n2. Description of Related Art\nA technique of calculating a maximum power point in solar battery arrays in total by connecting the solar battery arrays, which are connected in series to each other, to monitoring solar batteries corresponding to the number of solar battery arrays in parallel and detecting short-circuit currents of the monitoring solar batteries and open-circuit voltages of the solar battery arrays is known (for example, Japanese Patent Application Publication No. 61-281316 (JP 61-281316 A)).\nHowever, in the technique described in JP 61-281316 A, the timing of detecting the short-circuit currents of the monitoring solar batteries and the open-circuit voltages of the solar battery arrays is not described. Accordingly, when a solar battery module is mounted on a moving object, the maximum power point calculating timing cannot be appropriately set and thus there is a possibility that an increase in computational load or a delay in the maximum power point calculating timing will occur."}
{"text": "Power management is used in many devices and systems to improve power efficiency, helping to reduce power consumption and/or heat dissipation. For battery-powered mobile devices and systems, power management can help extend operation.\nSome platform-level power management may place a processor and/or a chipset into a lower power state which can impact input/output (I/O) performance. As one example where a platform supports bus mastering, an I/O device may initiate a data transfer to system memory when a processor and chipset are in a lower power state. The processor and chipset would have to return to a normal operating state to complete the data transfer and therefore introduce delay which reduces performance.\nThe figures of the drawings are not necessarily drawn to scale."}
{"text": "Particulate matter may be present in the exhaust of some engines, diesel engines for example. Vehicle tailpipe emissions may be reduced by installing particulate filters and particulate matter sensors in vehicle exhaust systems. One particular type of particulate sensor senses particulate matter by electrodes separated by a dielectric gap. When particulate matter bridges the gap between electrodes, the resistance and conductance of the particulate matter sensor is changed.\nThe above-mentioned sensor may under some conditions experience degradation. For example, a wire connecting the particulate matter sensor to an engine controller may become disconnected or holes in a protective shield may become partially plugged. In such circumstances it may be desirable to determine whether or not the particulate matter sensor is operating as desired. However, it may be difficult to establish sensor degradation because the particulate sensor exhibits a high resistance when there is little particulate matter between the positive and negative electrodes. Thus, the particulate sensor output may appear to look like an open circuit. Therefore, it may be difficult to distinguish a properly operating particulate sensor from an open circuit. One embodiment of the present description includes a method for detecting a degraded particulate matter sensor, comprising: starting an engine when a temperature of said engine is less than a threshold temperature; and indicating degradation of said particulate matter sensor in response to an output of said particulate matter sensor being below a threshold, said output related to water vapor in exhaust gases generated by said engine during said starting of said engine.\nThus, by recognizing that the output of a particulate matter sensor may be affected by water vapor contained in engine exhaust gases, the inventors herein have developed a way to diagnose a resistive particulate sensor that exhibits characteristics of an open circuit when the particulate sensor is operating properly in the presence of little particulate matter. For example, the particulate sensor can be monitored during an engine start when water in the exhaust may cause a change in the state of the output of the particulate sensor. If the output of the particulate sensor does not increase above a threshold level, it may be judged that the particulate matter sensor is degraded.\nThe present description may provide several advantages. Specifically, the approach may improve engine emissions by allowing degraded sensors to be replaced before engine emissions degrade. Further, special hardware for monitoring or stimulating the particulate sensor may not be required because the method may be performed with the hardware that controls the particulate heater and senses the particulate sensor output. As a result, system cost may be reduced.\nThe above advantages and other advantages, and features of the present description will be readily apparent from the following Detailed Description when taken alone or in connection with the accompanying drawings.\nIt should be understood that the summary above is provided to introduce in simplified form a selection of concepts that are further described in the detailed description. It is not meant to identify key or essential features of the claimed subject matter, the scope of which is defined uniquely by the claims that follow the detailed description. Furthermore, the claimed subject matter is not limited to implementations that solve any disadvantages noted above or in any part of this disclosure."}
{"text": "The prostaglandins are a group of hormone-like substances which may be viewed as derivatives of prostanoic acid. Several prostaglandins are found widely distributed in mammalian tissue and have been isolated from this source. These prostaglandins have been shown to process a variety of biological properties such as bronchodilation, the ability to reduce gastric secretion, to modify muscle tone, as well as the ability to raise or lower blood pressure.\nVarious derivatives of prostaglandins have also been synthesized and reported. 9,15-Dihydroxy prost-13-enoic acid and methods of synthesis thereof are disclosed in U.S. Pat. Nos. 3,432,541 and 3,455,922. 9-Oxo-15-hydroxy-15-methyl-prostanoic acid, 15-oxo-9-hydroxy-prostanoic acid, and 9,15-dioxo-prostanoic acid are disclosed in U.S. Pat. No. 3,671,570.\nThe present invention concerns new unsaturated 15-polycarbon(lower)alkyl, vinyl and benzyl 9-oxo-15-hydroxy-prostanoic acids and new intermediates thereto."}
{"text": "1. Field of the Invention\nThis invention relates to valve systems of the type adapted to allow the introduction of a surgical instrument into a patient's body. In particular, the invention is applicable to a cannula assembly wherein a cannula housing includes the valve assembly and the cannula is intended for insertion into a patient's body to sealingly accommodate an instrument inserted through the cannula and valve.\n2. Background of the Prior Art\nIn laparoscopic procedures surgery is performed in the interior of the abdomen through a small incision; in endoscopic procedures surgery is performed in any hollow viscus of the body through narrow tubes or cannula inserted through a small entrance incision in the skin. Laparoscopic and endoscopic procedures generally require that any instrumentation inserted into the body be sealed, i.e. provisions must be made to ensure that gases do not enter or exit the body through the incision as, for example, in surgical procedures in which the surgical region is isufflated. Moreover, laparoscopic and endoscopic procedures often require the surgeon to act on organs, tissues, and vessels far removed from the incision, thereby requiring that any instruments used in such procedures be relatively long and narrow.\nFor such procedures, the introduction of a tube into certain anatomical cavities such as the abdominal cavity is usually accomplished by use of a trocar assembly comprised of a cannula assembly and an obturator. The cannula assembly includes a cannula tube attached to a valve assembly which is adapted to maintain a seal across the opening of the cannula assembly. Since the cannula tube is in direct communication with the internal portion of the valve assembly, insertion of the cannula tube into an opening in the patient's body so as to reach the inner abdominal cavity must maintain a relatively gas-tight interface between the abdominal cavity and the outside atmosphere.\nSince surgical procedures in the abdominal cavity of the body require insufflating gases to raise the cavity wall away from vital organs, the procedure is usually initiated by use of a Verres needle through which a gas such as CO2 is introduced into the body cavity. Thereafter, the pointed obturator of the trocar assembly is inserted into the cannula assembly and used to puncture the abdominal cavity wall. The gas provides a slight pressure which raises the inner wall surface away from the vital organs thereby avoiding unnecessary contact with the organs by the instruments inserted into the cannula. Following removal of the obturator, laparoscopic or endoscopic surgical instruments may then be inserted through the cannula assembly to perform surgery within the abdominal cavity.\nIn view of the need to prevent leakage of the insufflation gas from the cavity, the cannula is typically provided with a valve assembly which permits introduction of surgical instruments to provide selective communication between the inner atmosphere of the cavity with the outside atmosphere. In this regard, there have been a number of attempts in the prior art to provide such a seal as part of the cannula assembly.\nOne form of cannula valve assembly includes a flapper valve which is pivotally mounted within the cannula assembly and is automatically opened by the obturator or other object when it is inserted into the proximal end of the cannula. Conventional flapper valves may also be manually opened by pivoting a lever on the exterior of the housing. See, e.g., U.S. Pat. No. 4,943,280 to Lander. Trumpet valves are also known.\nOther conventional cannula valve devices for accommodating surgical instruments include a single or plurality of flexible sealing members as shown, for example, in U.S. Pat. No. 4,655,752 to Honkanen et al., U.S. Pat. No. 4,909,798 to Fleischhacker, U.S. Pat. No. 4,673,393 to Suzuki et al., U.S. Pat. No. 4,610,665 to Matsumoto et al., and U.S. Pat. No. 4,869,717 to Adair.\nFurther, typical hemostasis valve devices are shown, for example, in U.S. Pat. No. 5,041,095 to Littrell, and U.S. Pat. No. 5,000,745 to Guest et al.,\nWhile attempts have been made to provide a valve assembly which maintains the integrity of the seal between the body cavity and the atmosphere outside the patient's body. Seal systems provided to date have failed to address the full range of surgeons' needs, especially when instruments varying in diameter are used. Specifically, sealing elements currently used may be damaged when an instrument, such as a pointed obturator is passed therethrough. Moreover, present seal systems have not provided adequate sealing about an instrument before and after an instrument is passed therethrough. Also, existing seal systems have failed to provide adequate sealing of a cannula, or a trocar assembly having a cannula which accommodates instruments of varying diameters. It is a further disadvantage of existing seal systems that adequate sealing is not provided in conjunction with a structure for holding a cannula in a desirable position in an incision with respect to a patient's body.\nIt would therefore be desirable to provide a valve assembly which addresses these shortcomings in the art by maintaining a substantially fluid tight seal between an internal portion of a patient's body and the outside atmosphere during insertion and manipulation of a surgical instrument into the patient's body. Such an assembly may further provide stabilization or lateral limitation of motion of an instrument passed therethrough. Also, the valve assembly may inhibit fluids from exiting with the instrument while being withdrawn, and the valve assembly may inhibit contact with sealing structure. It is further desirable to provide a valve assembly for use with a cannula or trocar assembly which provides substantial fluid and gas tight sealing before and after an instrument is passed therethrough. It would also be desirable to provide a cannula which maintains a predetermined position of a cannula or trocar assembly in an incision.\nThe present invention provides a valve assembly which may be incorporated into a cannula assembly or utilized in combination with any type of tubular member for providing access into the body of a patient while permitting introduction of instruments through the valve assembly into the body. The valve assembly includes a sealing gasket which provides a desirable seal about an instrument inserted through the valve assembly. The valve assembly may further provide stabilization of the cannula or limit lateral motion of the cannula when an instrument is passed therethrough. Also, the valve assembly may include more than one sealing element providing improved sealing qualities under varied conditions. At all times, the surgeon maintains control over the interface between the atmospheres within and without the patient's body. Moreover, the present invention makes it possible to introduce instruments of varying sizes into the body and insures the maintenance of a gas seal despite instrument manipulation therethrough."}
{"text": "1. Technical Field\nThe present invention relates to control devices, display devices and methods for controlling a display device.\n2. Related Art\nAmong electrophoretic type display devices, there is a type that performs each rewriting operation by applying voltages multiple times. Such a rewriting operation may be performed when it takes a relatively long time for display elements to change their display state (i.e., gradation). When such a rewriting operation is performed, unless one rewriting operation is completed (in other words, unless the time for multiple voltage applications has elapsed), the next rewriting operation cannot be started.\nJapanese Laid-open Patent Application JP-A-2009-251615 (Patent Document) describes a technology for rewriting an image at each of divided regions by a pipeline processing in a display device such as an electrophoretic display device or the like. By using such a method, a rewriting operation can be started on a region where rewriting is not performed, without depending on the rewriting state of other regions, such that the time required for rewriting may be made shorter, compared to the case where the entire image is rewritten at once.\nIncidentally, in an area of the electrophoretic display device where rewriting is performed a plurality of times, an afterimage of the image may remain due to operation failure of the electrophoretic particles. In order to erase such an afterimage, the entire target area is displayed in black or the entire target area is displayed in white, a so-called refreshing operation may be performed. However, according to the technology described in the Patent Document, partial refreshing is not described, but refreshing is performed over the entire display area. Therefore, while refreshing is performed, other rewriting operations are stopped."}
{"text": "The engine start system of the present invention generally relates to mechanical actuation of components used to start a turbojet engine in a more economical fashion than currently known systems.\nA one-shot start system for a turbojet engine is typically comprised of a pyrotechnic device called an ignitor activated by an initiator. Initiators (often called squibs) are generally electroexplosive devices utilizing an electrically heated bridgewire to ignite an explosive charge. The bridgewire is designed to ignite the charge only when a specific level of current is applied across it by a firing circuit. Since the initiator is electrically activated, the bridgewire and associated firing circuit must be carefully designed to prevent accidental activation from exposure to external electromagnetic signals produced by equipment such as conventional radars. Manufacturing an initiator wherein the bridgewire is electrically heated at a very precise firing current requires strict manufacturing and quality control processes that tend to increase the cost of the device.\nA bullet operates in an analogous fashion to the aforementioned ignitor. A sharp blow from the hammer of a gun mechanically fires an initiator, usually termed a primer. The primer then ignites a gunpowder charge which propels the bullet down a gun's barrel. The primer serves the same actuation function for a bullet as the bridgewire initiator does for an ignitor. However, unlike the initiator, the primer is a percussion-activated device. As such, it is much cheaper to produce than the initiator since it has no critical electronic components. Additionally, the primer used in a bullet is not susceptible to accidental activation from electromagnetic radiation as evidenced by the longstanding practice of using conventional bullets in combat aircraft applications. The use of a percussion primer initiator for turbojet engine start system will substantially reduce overall system cost.\nMost known turbojet engines utilize an electric fuel pump to meter fuel into an engine combustor. Engines using through-shaft fuel induction may utilize a low pressure fuel delivery system in conjunction with a pressurized fuel tank. This obviates the need for a high pressure fuel delivery device such as a pump. A low pressure fuel injector, for example an automotive injector, may be used in place of a metering fuel pump.\nConventional automotive injectors are typically operated by a solenoid actuated valve. These valves may be opened and closed by a pulse width modulated signal (PWM) which varies the valve actuation time by varying the signal pulse duration. (See for example U.S. Pat. No. 5,052,174, assigned to assignee of the present invention.) Fuel delivery systems using automotive type injectors have proven acceptable for use in turbojet engines when operating at pulse frequencies of 60 to 140 hertz. Automotive injectors are readily available, durable, and very cost effective when compared to metering pumps."}
{"text": "Metal, ferroelectrics, oxide, and silicon (MFOS) transistor ferroelectric memory devices have been proposed in the related applications. In order to provide a MFOS transistor memory device having desirable characteristics, the oxide must not react with nor diffuse into the ferroelectric and silicon substrate. On the other hand, ferroelectric thin films deposited on an oxide layer should have good ferroelectric properties for use in a memory transistor. Homogeneous ferroelectric thin films having desirable ferroelectric properties are difficult to deposit on gate oxides because of the interface mismatch between the high-k gate oxide and the ferroelectric materials. The mismatch results in random ferroelectric, such as PGO, thin films and coarser surface roughness. Seed layer processes have been developed for MFOS transistor ferroelectric memory applications to resolve the interface mismatch."}
{"text": "1. Field of the Invention\nThe field of the invention relates to the regulation of the growth of aquatic weeds in canals, rivers, ponds, lakes and impoundments.\n2. Description of the Prior Art\nThe problems of controlling or regulating the growth of organisms in aqueous systems are series and growing in severity. Submerged aquatic weeds, for example, cause major problems in water distribution and irrigation systems. The growth of such weeds in irrigation canals greatly reduces the conductivity and capacity of such systems with resulting substantial economic loss. Large sums are spent in the mechanical and other methods of removal of weed growths from irrigation canals, especially in the western parts of the United States. Because of the great difficulties involved in the mechanical removal of weeds and other undesired forms of aquatic life from irrigation canals, ponds, lakes, impoundments, etc., it has been proposed to utilize chemical control. Accordingly, various types of chemicals have been added to such bodies of water.\nHowever, with the growing emphasis on conservation, current efforts are directed toward regulating the growth, that is limiting or inhibiting the amount of growth accomplished by the naturally occurring submerged or floating aquatic weeds, without killing those weeds. This approach is being taken in order to continue to provide the natural environment for fish and other forms of marine life. A further reason is to avoid the masses of dead and rotting aquatic weeds which result when said weeds are killed by means of an aquatic herbicide, since the decomposition of the weeds decreases the amount of available oxygen present in the water. Such decaying matter, when it occurs in reservoirs and/or streams from which drinking water for cities is obtained, makes purification of the water much more difficult. Such decaying vegetation also gives off an unpleasant odor when it collects in a body of water. Thus, a control of the amount of growth rather than a destruction of the submerged aquatic weeds serves to overcome both pollution of the water and pollution of the air.\nIn the prior art, Krumkalns et al., U.S. Pat. No. 3,655,359 Apr. 11, 1972), teach and claim the use of substituted 3-pyridylmethanes for eliminating germinating weed grasses and broadleaf weeds selectively from crop plants such as corn, cotton, and soybeans, and their close relatives in the plant kingdom. I have found many of the compounds appearing in this reference useful for the regulation of the growth of aquatic weeds and they come within the scope of the compounds used in the present invention. There is no teaching in the reference that the compounds disclosed therein would act to control the growth of submerged or floating aquatic weeds.\nAlso in the prior art, Krumkalns et al., U.S. Pat. No. 3,744,988 (July 10, 1973), teach and claim the use of substituted 3-pyridylmethanes in a method for inhibiting sucker growth to tobacco plants. This patent is a division of U.S. Pat. No. 3,655,359, supra, and includes in its disclosure many of the same compounds disclosed in that patent. There is no teaching in the reference that the compounds disclosed therein would act to control the growth of submerged or floating aquatic weeds.\nVan Heyningen, U.S. Pat. No. 3,396,224 (Aug. 6, 1968), teaches and claims a method of controlling fungi pathogenic to plants by contacting the fungus-susceptible plant with a fungicidal amount of a 3-pyridylmethane derivative, mainly a 3-pyridinemethanol. I have found many of the compounds disclosed in this reference active as aquatic plant growth regulators, and they are included within the scope of the instant application. There is no teaching in this reference that the compounds disclosed would be active as aquatic growth regulators of submerged or floating aquatic weeds.\nVan Heyningen et al., U.S. Pat. No. 3,397,273 (Aug. 13, 1968), teach and claim a method for protecting plants from attack by phytopathogenic fungi by treating the plants with a fungicidally-effective amount of a 3-pyridylmethane. Some of the compounds disclosed by this reference have now been found by me to be active as aquatic growth regulators, and their use for that purpose is disclosed and claimed in the instant application. This reference makes no suggestion that the therein disclosed compounds would be active as growth regulators of submerged or floating aquatic weeds.\nKrumkalns, U.S. Pat. No. 3,335,148 (Aug. 8, 1967), discloses and claims the 9-(3-pyridyl) derivative of fluorene, 9-fluorenol, xanthene, 9-xanthenol, and the corresponding nonphytotoxic acid addition salts thereof. The compounds are alleged to posses antifungal and antibacterial activities. I have now found some of the compounds to be active as aquatic growth regulators, and such compounds are included for that use in the instant application. There is no teaching or suggestion in this reference that the compounds would be active as growth regulators of submerged or floating aquatic weeds.\nYet another reference, Krumkalns, U.S. Pat. No. 3,361,753 (Jan. 2, 1968), is directed to 9-(3-pyridyl)-thioxanthene and thioxanthol derivatives, active as plant antifungal agents and as antibacterial agents. There is no teaching that the compounds of this reference would be active as regulators of the growth of submerged or floating aquatic weeds.\nAnother reference is German Pat. No. 1,935,292, also identified by Derwent No. 04548S, which patent teaches and claims a means for controlling plant growth, that is restraining growth and influencing the habits of higher plants, influencing blossom and fruit formation, checking the growth of grass, and the like, using triarymethylimidazoles, -pyrazoles, and -triazoles, or their salts. One of the aryl groups is taught as pyridyl. The reference does not include use on aquatic weeds or plants.\nYet another reference is British Patent No. 1,274,578, also identified by Derwent No. 23143S. This reference teaches plant growth regulators containing N-benzylimidazoles, wherein one of the substituents is a pyridyl group. These compounds are alleged to be plant growth regulators capable of inhibiting or accelerating growth, flowering, and fruiting, according to the amount applied. Certain of the compounds are also alleged to be plant fungicides and bactericides."}
{"text": "Tie-downs for mounting objects on motor vehicle frames are well known. Such tie downs are often provided for fixed installation to the frame, or to the underside of a vehicle (such as a pick-up truck), for securing a load (such as a camper) to the vehicle frame. Generally, four such tie downs are installed on a vehicle for securing the camper to the pick-up truck's bed at or near all four corners. The currently used tie downs known to me generally include a cantilevered telescoping arm projecting horizontally from a secure position on the vehicle frame. Such prior art tie downs include a rigid support member projecting perpendicularly from an outer sleeve of a cantilevered telescoping arm to fixedly attach to a point on the underside of the vehicle to brace the tie down when the load is secured. Also, such tie downs include a chain having one end pliantly attached to the end of the telescoping arm opposite the frame mount and the other end free for fixedly attaching to the load, such as a camper, to secure it to the vehicle. The load is secured to the vehicle by tensioning the chain using, for example, a turnbuckle. Thus tensioned, the chain is a rigid member unable to relieve sudden loads, such as those caused by shifting of the camper relative to the truck bed. Such sudden load, when not relieved, may exceed the stress-bearing capabilities of the joints mounting the tie downs to the vehicle and may loosen the joints and/or damage either or both of the vehicle and the tie downs. While a spring-loaded turnbuckle is known, it is a separate assembly from the tie down and forms no part of the tie down.\nIt is a drawback that the currently utilized tie downs must be professionally installed by qualified mechanics, primarily since they are provided from the factory in a single embodiment or configuration which requires drilling of additional mounting holes in the vehicle frame. The drilling of such holes in or adjacent to highly loaded or specially hardened steel frame members requires equipment and a skill level ordinarily available only to the professional installer. Furthermore, the warranties provided by many vehicle manufacturers might be violated or voided by the addition of such holes in the vehicle frame members. Therefore, such considerations severely limit the practical application of the known tie downs, as well as effectively prevent their distribution by mail order.\nUndesirably, many tie downs also interfere with use of the vehicle even when the load is not present. This is because rigid support members projecting perpendicularly from the outer sleeve of the cantilevered telescoping arm places the tie down in a position beneath the vehicle in a manner that may interfere with steps, bars and other platforms attached to many vehicles for ease of entry and exit. Fixed position tie downs may also interfere with use of the vehicle off-road, because they are positioned below the vehicle, near the ground, and may encounter brush, stumps, rocks and other low lying obstacles."}
{"text": "1. The Field of the Invention\nThe present invention relates generally to techniques for measuring levels of chemical compounds found in biological tissue. More specifically, the invention relates to a method and apparatus for the noninvasive detection and measurement of levels of carotenoids and related chemical substances in biological tissue, which can be used as a diagnostic aid in assessing antioxidant status and detecting malignancy diseases or risk thereof.\n2. The Relevant Technology\nCarotenoids are plant pigments available from the diet which have important functions in the human body. The role of carotenoids in human health is a rapidly expanding area of research. Much carotenoid research has focused on their role as precursors to retinoids or vitamin A, but current research is also being conducted on other functions of carotenoids. These include antioxidant activities, modulation of the immune response, cell-to-cell communication and gap junction modulation.\nIt has been demonstrated that carotenoids offer some degree of biologic protection against the formation of malignancies in various tissues. For example, carotenoids have been shown in animal models to prevent carcinoma formation in tissues such as skin, salivary gland, mammary gland, liver, and colon. In addition, low levels of carotenoids and related substances such as retinoids have been assessed as high risk factors for malignant lesions. For example, having low levels of the carotenoid lycopene has been associated with prostate and cervical cancer; the carotenoids lutein, zeaxanthin, alpha-carotene, and beta-carotene with lung cancer; and beta-carotene with oral cancer. Therefore, quantitatively measuring the chemical concentrations of these carotenoids, retinoids and other related substances provides an indicator of the risk or presence of cancer.\nThe most common cancer in the United States is skin cancer. Despite attempts at patient education, skin cancer rates continue to rise. Methods to provide detection of the levels of chemicals which are associated with skin related malignancies are of great assistance to physicians and medical personnel in the early diagnosis and treatment of skin cancer.\nIt has been theorized that carotenoids in the skin provide biologic protection from cutaneous malignancy. Most findings, however, have been somewhat compromised by the fact that concentrations of carotenoids in skin and skin malignancies were never measured directly, with data on levels of carotenoids in patients being derived only indirectly from blood plasma.\nPrior methods used to detect the presence of chemicals associated with skin cancer have mainly been through the analysis of tissues obtained by biopsies or other invasive procedures. The standard method presently used for measuring carotenoids is through high-performance liquid chromatography (HPLC) techniques. Such techniques require that large amounts of tissue sample be removed from the patient for subsequent analysis and processing, which typically takes at least twenty four hours to complete. In the course of these types of analyses, the tissue is damaged, if not completely destroyed. Therefore, a noninvasive and more rapid technique for measurement is preferred.\nA noninvasive method for the measurement of carotenoid levels in the macular tissue of the eye is described in U.S. Pat. No. 5,873,831, the disclosure of which is herein incorporated by reference, in which levels of carotenoids and related substances are measured by a technique known as Raman spectroscopy. This is a technique which can identify the presence and concentration (provided proper calibration is performed) of certain chemical compounds. In this technique, nearly monochromatic light is incident upon the sample to be measured, and the inelastically scattered light which is of a different frequency than the incident light is detected and measured. The frequency shift between the incident and scattered light is known as the Raman shift, and the shift corresponds to an energy which is the \"fingerprint\" of the vibrational or rotational energy state of certain molecules. Typically, a molecule exhibits several characteristic Raman active vibrational or rotational energy states, and the measurement of the molecule's Raman spectrum thus provides a fingerprint of the molecule, i.e., it provides a molecule-specific series of spectrally sharp vibration or rotation peaks. The intensity of the Raman scattered light corresponds directly to the concentration of the molecule(s) of interest.\nOne difficulty associated with Raman spectroscopy is the very low signal intensity which is inherent to Raman scattered light. It is well known that the scattered light intensity scales with the frequency raised to the fourth power. The weak Raman signal must be distinguished from Rayleigh scattered light, which is elastically scattered light of the same frequency as the incident light and which constitutes a much greater fraction of the total scattered light. The Raman signal can be separated from Rayleigh scattered light through the use of filters, gratings, or other wavelength separation devices; however, this can have the effect of further weakening the measured Raman signal through the additional attenuation which can occur when the light passes through a wavelength separation device. In practice, the Raman scattered light is extremely difficult to detect. One might attempt to increase the Raman signal by increasing the incident laser power on the tissue sample, but this can cause burning or degradation of the sample.\nIn order to overcome some of these difficulties, a technique known as resonance Raman spectroscopy has been used, as described in U.S. Pat. No. 5,873,831, referenced hereinabove. Such a technique is also described in U.S. Pat. No. 4,832,483, the disclosure of which is herein incorporated by reference. In resonance Raman spectroscopy, the incident illumination utilized has a frequency which corresponds to the resonance frequency corresponding to electronic energy transitions of the molecules of interest. This has the effect of strongly enhancing the Raman output signal without using a higher intensity input signal, thereby avoiding damage to the sample which can be caused by laser burning. Also, these resonance Raman signals have much higher intensity than off-resonance Raman signals which are virtually invisible. Therefore, in resonance Raman spectroscopy only those Raman signals which belong to the species of interest are obtained.\nIn the above referenced U.S. Pat. No. 5,873,831, the resonance Raman technique is used to measure the levels of the carotenoids lutein and zeaxanthin, two chemicals which are associated with healthy macular tissue of the human eye. The above referenced U.S. Pat. No. 4,832,483 uses resonance Raman spectroscopy to measure certain carotenoids in blood plasma, and suggests the use of the ratios of the intensities of the Raman spectral peaks as a method of indicating the presence of various malignancy diseases.\nYet another difficulty associated with Raman measurements is that the substances of interest in the skin not only scatter incident light, but can absorb and subsequently fluoresce with substantial intensity. This fluorescence often comprises a very strong, broad signal which tends to \"drown out\" or overwhelm the Raman spectral peaks, thereby making identification and quantification of the substances of interest practically impossible.\nFluorescence spectroscopy is itself another technique which can be used to measure amounts of chemical compounds in biological tissue. For example, U.S. Pat. No. 5,697,373 discloses use of fluorescence and/or Raman spectroscopy to detect tissue abnormality in the cervix. The disadvantage of fluorescence measurements is that since many different molecules fluoresce in broad bands of frequencies, such measurements cannot be used to conclusively identify the presence or concentration of a particular substance.\nIt would therefore be a significant advance to provide a method and apparatus for the safe, noninvasive, rapid, accurate, and specific measurement of the levels of carotenoids and other similar chemical compounds which are present in varying degrees in biological tissues, and to use this information to aid in the assessment of cancer risk or disease risk in all types of biological tissue."}
{"text": "1. Field of the Invention\nThe subject invention relates to containers, such as for example corrugated containers, paperboard containers, etc. having tear lines and blanks for making such containers. It should be understood by those of ordinary skill in the art that the word container is being used broadly to include, but not be limited to boxes, cartons, bags etc. and not limited in anyway to the material used.\n2. Description of the Prior Art\nToday, many types of containers have been developed for the attractive packaging of articles, such as food products, drug items, cosmetics, and the like. Conventionally, the containers are formed from a one-piece blank of material which is cut, scored and folded to produce a generally rectangularly shaped box.\nThe art is replete with various other designs of cartons, which are disclosed in U.S. Pat. No. D436,859 to Botsford et al.; U.S. Pat. No. 5,881,884 to Podosek; U.S. Pat. No. 5,979,749 to Bozich; U.S. Pat. No. 6,073,833 to Desrosiers et al.; U.S. Pat. No. 6,102,277 to Krapohl, Sr.; U.S. Pat. No. 6,435,351 to Gibb; U.S. Pat. No. 6,752,262 to Boriani et al.; U.S. Pat. No. 6,766,941 to Tokarski; and U.S. Pat. No. 6,935,557 to Aubry et al.\nThe U.S. Pat. No. 5,979,749 to Bozich, for example, teaches a carton that includes a front panel, a back panel, and side panels. The carton includes a top panel and a bottom panel. The carton includes a tear line extending about the periphery of the carton. The tear line extends into a gripper used by a consumer (not shown) to tear the tear line to expose the carton to an open top as a top cover is removed from the carton. A display section is defined in the front panel and can be separated from the front panel by a perforation or tear line. All items contained in the carton are exposed to and are accessible to the consumer when the top cover is removed from the carton. The items are only available for display purposes when the display section is severed from the front panel.\nThe U.S. Pat. No. 6,102,277 to Krapohl, Sr., for example, teaches a carton having a front panel, a back panel, and side panels. The carton includes a top panel and a bottom panel. The carton includes a tear line extending about periphery of the carton. The tear line extends into a gripper used by a consumer to tear the tear line to expose the carton to a first opening as the bottom of the carton is removed. The top panel is defined by a pair of flaps with each flap being integral with and extending from the front panel and the back panel, respectively. The flaps are mechanically interconnected with one another and present a second opening. The carton is designed for holding a bag containing perishable goods. One of the most common problems associated with the prior art patents is a configuration of the perforations or slits along the tear line, which results in the corners, defined between the front, back, and side panels, being damaged and deformed.\nThere is a constant need in the area of the packaging industry for an improved carton adaptable for dual dispensing of a product contained therein at a user's preference and an improved configuration of the perforations or slits along the tear line."}
{"text": "This application is related to the below listed co-pending patent applications which are filed on even date herewith, are assigned to the same assignee, and are incorporated herein in their entirety by these references:\nAn application entitled xe2x80x9cMethod and Apparatus For Interactively Selecting Display Parameters For An Avionics Flight Displayxe2x80x9d by Sarah Barber, Norm W. Arons, and George W. Palmer;\nAn application entitled xe2x80x9cMethod and Apparatus For Interactively Selecting, Controlling and Displaying Parameters For An Avionics Radio Tuning Unitxe2x80x9d by George W. Palmer, Claude Eyssautier, and Matt Smith;\nAn application entitled xe2x80x9cMethod and Apparatus For Interactively Displaying A Route Window For A Flight Management Systemxe2x80x9d by Gary L. Owen, Sarah Barber, and George W. Palmer; and\nAn application entitled xe2x80x9cMethod and Apparatus For Graphically Inserting Waypoints For A Flight Management Systemxe2x80x9d by Martin Pauly.\nThe present invention generally relates to avionics, and more particularly relates to avionics display systems, and even more particularly relates to multi-functional electronic flight display systems.\nIn the past, designers of avionics displays and flight deck electronic navigation, control and communication systems have endeavored to achieve a reduction in pilot workload. One area of concern has been the control of new large-sized multi-functional displays (MFDs), which typically requires a significant amount of xe2x80x9cheads-downxe2x80x9d time. This xe2x80x9cheads-downxe2x80x9d time occurs when the pilot is neither looking at the primary flight displays, nor out the wind screen, but instead is focused upon a task in an oblique direction, such as when manually controlling the information content of these MFDs via a separate controller, which may or may not be co-located with the MFDs and may have an integrated keypad or tuning knob console. Another approach has been proposed in which MFDs are segmented into smaller viewing areas and are controlled via a single pop-up, cursor-driven menu in the center of the MFD. The pilot can select from several predetermined options. Depending on the predetermined size of each instrument function to be displayed, earlier selected functions may be either removed from the view, and/or shifted to a new position on the MFD. While these alternate means of controlling MFDs have clear advantages, they also have significant drawbacks.\nDuring times of moderate and severe turbulence, otherwise very simple tasks can become too difficult and too time consuming. For example, the step of requiring a pilot to look for or mentally keep track of the shifting location of each previously displayed instrument function is undesirable. Requiring the pilot\"\"s attention to be diverted to a side-mounted display to operate a separate controller may also be undesirable in certain circumstances. In general, during take-off and approach, the workload on a pilot can already be extreme, leaving the pilot with little or no time to spare. Any additional effort, at an already busy time, is quite undesirable.\nConsequently, there exists a need for improved methods and apparatuses for effecting the operation of electronic flight instruments in an aircraft which overcome some of these shortcomings.\nIt is an object of the present invention to provide an improved avionics electronic flight display system.\nIt is a feature of the present invention to include a segmented forward-mounted multi-functional display with enhanced menu display capabilities, including localized pop-up menus which appear in each segment upon clicking in such segment.\nIt is another feature of the present invention to include a graphical user interface which allows for automatic vertical expansion/compression of a segment of an electronic flight instrument display.\nIt is yet another feature of the present invention to include a manual expansion/compression feature.\nIt is still yet another feature of the present invention to include automatic presentation of a one-click toggle feature, for linking to a hidden page, when an expanded display segment is compressed.\nIt is even another feature of the present invention to include a feature to prohibit simultaneously displaying the same instrument in two or more display segments.\nIt is an advantage of the present invention to reduce pilot workload.\nIt is another advantage of the present invention to enhance pilot situational awareness.\nThe present invention is an apparatus and method for selecting, controlling and displaying parameters for an avionics electronic flight display system, which is designed to satisfy the aforementioned needs, provide the previously stated objects, include the above-listed features, and achieve the already articulated advantages. The present invention is carried out in a xe2x80x9cheads down-lessxe2x80x9d manner in a sense that the undesirable requirement for the pilot\"\"s attention to be diverted from the primary flight displays is reduced.\nAccordingly, the present invention is an avionics multi-functional display with enhanced menuing capabilities and operating in innovative ways so as to reduce pilot xe2x80x9cheads-downxe2x80x9d time."}
{"text": "1. Field of the Invention\nThis invention relates to hardware and software installation and maintenance, and more particularly to software programs for diagnosing product issues.\n2. Description of the Related Art\nComputer networks and other products may have multiple components. When a product issue affects one component, the product issue may eventually affect the other similar components on the products in a comparable way. For example, if an installer installed several similar components on multiple products, the same error may have been made in each installation. Once the error is detected on one component, it may need to be remedied on similar components of other products. In addition, many product issues with components may not be detected until later if product issue symptoms are delayed. In addition, product issues discovered on one product may affect other similar products over the course of the product's lifetime.\nExpert repairmen and on-site expert personnel may fix many product issues. Repair manuals may be consulted to aid with unfamiliar product issues. In addition, experts may watch or consult other experts to find out how to fix a product issue they are unfamiliar with. However, the spread of knowledge from expert to expert may be slow and incomplete. In many repair instances, the repairs for similar product issues may not be uniform and therefore, the results of these repairs may be unreliable. Furthermore, product issue solutions may change with time. Previously repaired products may need to be repaired again or inconsistent repairs across products may affect the product's reliability."}
{"text": "Multiple input multiple output (MIMO) increases the efficiency of data transmission and reception using multiple transmit antennas and multiple receive antennas instead of a single transmission antenna and a single reception antenna. A receiver receives data through multiple paths when multiple antennas are used, whereas the receiver receives data through a single antenna path when a single antenna is used. Accordingly, MIMO can increase a data transmission rate and throughput and improve coverage.\nA single cell MIMO scheme can be classified into a single user-MIMO (SU-MIMO) scheme for receiving a downlink signal by a single UE in one cell and a multi user-MIMO (MU-MIMO) scheme for receiving a downlink signal by two or more UEs.\nChannel estimation refers to a procedure for compensating for signal distortion due to fading to restore a reception signal. Here, the fading refers to sudden fluctuation in signal intensity due to multipath-time delay in a wireless communication system environment. For channel estimation, a reference signal (RS) known to both a transmitter and a receiver is required. In addition, the RS can be referred to as a RS or a pilot signal according to applied standard.\nA downlink RS is a pilot signal for coherent demodulation for a physical downlink shared channel (PDSCH), a physical control format indicator channel (PCFICH), a physical hybrid indicator channel (PHICH), a physical downlink control channel (PDCCH), etc. A downlink RS includes a common RS (CRS) shared by all user equipments (UEs) in a cell and a dedicated RS (DRS) for a specific UE. For a system (e.g., a system having extended antenna configuration LTE-A standard for supporting 8 transmission antennas) compared with a conventional communication system (e.g., a system according to LTE release-8 or 9) for supporting 4 transmission antennas, DRS based data demodulation has been considered for effectively managing RSs and supporting a developed transmission scheme. That is, for supporting data transmission through extended antennas, DRS for two or more layers can be defined. DRS is pre-coded by the same pre-coder as a pre-coder for data and thus a receiver can easily estimate channel information for data demodulation without separate precoding information.\nA downlink receiver can acquire pre-coded channel information for extended antenna configuration through DRS but requires a separate RS other than DRS in order to non-pre-coded channel information. Accordingly, a receiver of a system according to LTE-A standard can define a RS for acquisition of channel state information (CSI), that is, CSI-RS."}
{"text": "The present invention relates to a system for automatically controlling the idling speed of an internal combustion engine, comprising a valve for supplying an adjustable quantity of additional air and generally set so as to choke a duct connecting zones up- and downstream from the throttle valve controlled by the accelerator.\nOn known automatic idling control systems, said valve consists of an electrovalve, the setting of which is controlled by a signal as a function of the difference between required and actually detected engine speed, for the purpose of maintaining engine speed constantly within a given range under varying operating conditions.\nA major drawback on such known control systems is that they fail to provide for fast adjustment response, the disadvantage of which is particularly felt when applied to electronic injection systems providing for highly accurate, high-speed overall control of the engine via an electronic control system which, depending on signals from various sensors (mainly engine speed/stroke and air intake pressure/temperature sensors), determines, for example, air density inside the manifold, and engine speed, and calculates, via interpolation on respective memorized maps, the stroke and timing for injecting fuel into the injectors, as well as the spark lead."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device using a chip on film (COF) operating with a low driving voltage.\n2. Background of the Related Art\nLately, display devices having wide image areas are the trend of the times. The display devices have much more data and scan lines in order to improve image quality.\nHowever, if the number of the scan lines is increased to improve the image quality of the display device having a simple matrix structure, there are many limitations on driving the display device.\nFIG. 1 illustrates a diagram of a display device using a COF according to a related art.\nReferring to FIG. 1, data lines are connected straight to a chip on film (hereinafter abbreviated COF) from a display panel. And, scan lines, which are arranged on left and right sides of the display panel respectively, extend from the display panel with forming a curve so as to be coupled with the COF.\nAs mentioned in the above explanation, the scan lines forming curves from both sides of the display panel extend longer than the data lines. Such a structure results in increasing line resistance of the scan lines and further delays a signal response.\nIn display devices such as EL (electro-luminescence), LED (light emitting diode), FED (field emission display) and the like, the line resistance of the scan lines has a close relation to the driving voltage of the display devices.\nNamely, if the line resistance of the scan lines increases, the driving voltage thereof increases as well. Thus, the line resistance of the display device having a low driving voltage is a major issue.\nAccordingly, the present invention is directed to a display device using a COF that substantially obviates one or more problems due to limitations and disadvantages of the related art.\nAn object of the present invention is to provide a display device using a COF enabling to reduce line resistance of scan and data lines in a display device.\nAdditional advantages, objects, and features of the invention will be set forth in part in the description which follows and in part will become apparent to those having ordinary skill in the art upon examination of the following or may be learned from practice of the invention. The objectives and other advantages of the invention may be realized and attained by the structure particularly pointed out in the written description and claims hereof as well as the appended drawings.\nTo achieve these objects and other advantages and in accordance with the purpose of the invention, as embodied and broadly described herein, a display device using a COF according to the present invention includes a display panel, a COF having a connection part on which the display channel is placed and a chip part on which a chip is placed, data and scan lines extending from the display panel, scan line connection lines on the connection part, the scan line connection lines connecting the scan lines electrically to the COF, and data line connection lines on the connection part, the data line connection lines connecting the data lines electrically to the COF, wherein the scan lines are directly connected to the chip on a first line in parallel with an extending line from the COF through the scan line connection lines formed on the COF at both sides of the display panel.\nPreferably, the scan line connection lines are straightly connected to the scan lines at both sides of the display panel corresponding to the first line and the data line connection lines are connected to the data lines extending from both sides of the display panel corresponding to the second line. Also, a substrate material of the COF is one of glass and plastic.\nIn another aspect of the present invention, a display device using a COF includes a display panel, the COF having a connection part on which the display channel is placed and a chip part on which a chip is placed, data and scan lines extending from the display panel, scan line connection lines on the connection part, the scan line connection lines connecting the scan lines electrically to the COF, and data line connection lines on the connection part, the data line connection lines connecting the data lines electrically to the COF, wherein the scan lines are straightly connected to the COF through the scan line connection lines at both sides of the display panel on a first line in parallel with a direction extending from the COF, a portion of the data lines extends toward the chip from both sides of the display panel corresponding to a second line in a direction different from the extending direction of the COF, and the other portion of the data lines extends toward the chip from both sides of the display panel corresponding to a second line in a direction different from the extending line of the COF and is connected to the COF through the data line connection lines.\nIn a further aspect of the present invention, a display device using a COF includes a display panel, a COF having a connection part on which the display channel is placed and a chip part on which a chip is placed, data and scan lines extending from the display panel, scan line connection lines on the connection part, the scan line connection lines connecting the scan lines electrically to the COF, and data line connection lines on the connection part, the data line connection lines connecting the data lines electrically to the COF, wherein the scan lines are straightly connected to the COF toward the chip through the scan line connection lines at both sides of the display panel on a first line in parallel with an extending direction of the COF, and the data lines are connected to the COF through the data line connection lines from both sides of the display panel on a second line in a direction different from the extending direction of the COF.\nPreferably, the scan line connection lines extend straightly on the first line and the data line connection lines form a curve with a predetermined angle so as to extend from the first line to the second line. The predetermined angle is 90xc2x0.\nIt is to be understood that both the foregoing general description and the following detailed description of the present invention are exemplary and explanatory and are intended to provide further explanation of the invention as claimed."}
{"text": "As electronic and digital content use in enterprise settings and/or other organizational settings has become the preferred mechanism for project, task, and work flow management, so has the need for streamlined collaboration and sharing of digital content and documents. In such collaboration environments, multiple users share, access, and otherwise perform actions or tasks on content and files in shared workspaces.\nWhen a user performs an action on a file in a collaboration environment, a corresponding job can be scheduled. For example, in response to a file being uploaded, the file might responsively be scanned. Current automation architectures for collaboration environments provide a mechanism to kick off the jobs at the front-end (e.g., at the web or application servers). Unfortunately, these current architectures are not easily scalable and do not provide for customizations of the jobs to be performed responsive to particular actions in a distributed computing environment."}
{"text": "The present invention relates to a method for detecting an aberrant animal-derived prion gene. The invention further relates to a prion gene modified mouse which exhibits heart anomalies, and to a method for detecting drugs which reduce anomalies in an electrocardiogram of said mouse.\nTransmissible spongiform encephalopathy is a neurodegenerative disease observed in many kinds of mammal. These diseases are prion diseases caused by anomalies in prion proteins (PrP). Among these, kuru, Creutzfeldt-Jacob disease (CJD), Gerstmann-Straussler-Scheinker desease (GSS) in humans; and scrapie and bovine spongiform encephalopathy in ruminant livestocks such as sheep and cattle are known, and are a serious problems even today. A cause of the epidemic of prion diseases such as BSE in animals is thought to be livestock feed from sheep infected with scrapie, however, the causal link between the feed and incidence has not been confirmed.\nThe British medical journal xe2x80x9cLancetxe2x80x9d issued on Apr. 6, 1996 reported 10 cases of human transmissible spongiform encephalopathy, a new type of Creutzfeldt-Jacob disease (CJD) in the United Kingdom. This disease has been on the increase having killed 48 persons in the United Kingdom and 2 persons in France; and CJD has been under surveillance all over the world including Japan. The World Health Organization (WHO) has promoted research and development using transgenic animals. They are considered to be useful for detecting diseases sensitively in early stages and developing removal methods of pathogens. Owing to this research, the causal link between mad cow disease and CJD is gradually being elucidated.\nAt present, in Prusiner\"\"s report (S. Prusiner et al. PNAS 96:15137-15192, 1999) a bovine prion gene modified mouse is presented; and each pathogen, BSE, scrapie and CJD is said to appear at around 200 days in this transgenic mouse. Accordingly, the transgenic mouse indicates the high possibility that these three pathogens are identical. On the other hand, in Collinge\"\"s report (J. Collinge et al. PNAS, 91:9936-9940, 1994), a human prion gene modified mouse was presented; and CJD is said to appear at around 220 days in this mouse. However, a mouse in which the incubation period has been shortened using various animal-derived prion genes, is not yet known.\nIt is of great interest to note the fact that a wild ruminant such as kudu and oryx has a shorter incubation period compared with that of livestock ruminants such as sheep and cattle; and their progress after incidence is rapid. For example, it has been reported that generally the incubation period in sheep is 36-48 months and incidence is at age 36-48 months, and in cattle the incubation period is 36-72 months and the incidence is at age 60-80 months, but on the other hand, in Oryx demmah the incubation period is 21 months and incidence is at age 30 months. This indicates that the difference of transmissibility of spongiform encephalopathy and the incubation period might depend on the difference of amino acid sequence of each PrP. Up to now, the gene encoding a prion protein has been identified in many mammals such as human, sheep, cattle, mouse, hamster, etc; and it is known that their genes show not less than 90% amino acid homology to each other and they are highly conserved (N. Oesch et al., Cell, 40, 735 (1985); D. Westaway et al., Proc. Nat. Acad. Sci. USA 91: 6418-6422 (1994); D. Westaway et al., Genes and Develop. 8:959-969 (1994); N. Hunter et al., J. Gen. Virol. 74 1025-1031 (1993)).\nNow the system for rapidly detecting the scrapie pathogen is required to be established. Human or bovine prion gene modified transgenic mice have not yet been provided sufficiently in the world; therefore under present circumstances, the pathogen is detected by inoculating the animal tissues of interest in the brain of an ordinary CD-1 mouse (CLEA JAPAN). In the case of inoculating sheep brain tissue with this method, at least 300 days are required for the detection. Establishing a method for detecting the aberrant prion gene within a shortened incubation period before prion disease incidence may overcome the disadvantage of detection by the present bioassay which has high sensitivity but lacks rapidness.\nOn the other hand, a heart disease model in mouse is rare. Although the transgenic mouse with familial amyloid polyneuropathy has been reported, there is no report of a transgenic mouse as a heart specific disease model. The prion gene modified mouse with heart specific disease may also be useful for developing pharmaceuticals for heart diseases and studies on life-style related diseases.\nThe object of the present invention is to provide a method for detecting an aberrant animal-derived prion gene.\nThe present inventors have studied extensively and intensively to solve the above problems and have now found that a prion gene modified mouse exhibits heart disease when a foreign prion gene is aberrant, thereby completing the invention.\nFurthermore, in the past, there was no report which confirms abnormal electrocardiogram wave profiles and abnormal tissues in the heart of a mouse which overexpresses mouse normal prion genes (m-Prnp) (Prion and Prion Disease, Jun Tateishi. KYORITSU SHUPPAN, p23 (1998)).\nThe present invention relates to a method for detecting an aberrant animal-derived prion gene in which the prion gene is determined to be aberrant when the prion gene modified mouse exhibits heart disease.\nThe invention further relates to a prion gene modified mouse which exhibits heart disease.\nStill further, the invention relates to a method for detecting drugs which reduce abnormal waves in an electrocardiogram of the prion gene modified mouse.\nThe invention will be described in detail as follows."}
{"text": "1. Field of Invention\nThe invention is directed to downhole milling tools utilized in oil and gas wells to abrade, cut, or mill an object within the well and, in particular, to downhole cutting tools having a blade that is retracted during run-in and extended radially outward for cutting, abrading, or milling.\n2. Description of Art\nIn the drilling, completion, and workover of oil and gas wells, it is common to perform work downhole in the wellbore with a tool that has some sort of cutting profile interfacing with a downhole structure. Examples would be milling a downhole metal object with a milling tool or cutting through a tubular with a cutting or milling tool. To facilitate these operations, cutting elements are disposed on the downhole cutting tool. In one type of milling tool, the cutting elements are disposed on blades that can be disposed in a retracted position and an extended position. In certain of these embodiments, the blades are extended by increasing the pressure across the tool. Upon reduction of the pressure, such as after the milling operation has been completed, the blades are moved back to their retracted position so that the tool can be retrieved from the well."}
{"text": "1. Field of the Invention\nThe disclosed concept relates to a switchgear and, more particularly, to a switchgear wherein an enclosure housing assembly includes an extendable conductor assembly and a ground bus bar assembly that does not extend over the path of travel of a movable carriage assembly.\n2. Background Information\nElectrical switching apparatus, such as but not limited to circuit breakers, are often disposed in an enclosure or housing assembly. The electrical switching apparatus is disposed on, or is incorporated with, a movable carriage assembly. When the electrical switching apparatus need repair or other maintenance, the electrical switching apparatus is drawn, at least partially, from the enclosure. The removal of the electrical switching apparatus from the housing assembly may, or may not, decouple the electrical switching apparatus from the power source (hereinafter “line”) and/or the load. For the safety of the operators, a movable carriage assembly ground connection, i.e. a grounded conductor, is required by various regulations and safety procedures. For example, the frame of the circuit breaker is, typically, an electrically conductive material which needs to be coupled to ground whenever secondary wiring connections mate and/or become energized with electrical power\nThe ground connection allows the electrical switching apparatus to be in electrical communication with the ground continuously from the fully inserted/connected position until the electrical switching apparatus is drawn out a safe distance from the housing assembly. The purpose of this grounding feature is to protect operators from injury as a result of an electrical issue in the secondary control voltage system, which is typically less than 250 volts. The ground connection is also used with Ground and Test Devices that connect the primary voltage terminals, which typically can carry up to 15,000 volts, to ground so that the operators are protected while working on the conductor buses.\nTo allow for such a continuous ground connection, the housing assembly includes a large conductive bar (hereinafter “ground bar”), such as, but not limited to a copper bar, secured to the housing assembly. Further, the electrical switching apparatus includes a sliding electrical coupling that is structured to be slidably coupled to, and in electrical communication with, the ground bar. That is, the ground bar extends over the path of travel of the electrical switching apparatus so that as the electrical switching apparatus moves out of, or into, the housing assembly, the sliding electrical coupling moves along the ground bar maintaining the electrical switching apparatus in a grounded state.\nBecause advancing technology continually reduces the overall size of an electrical switching apparatus and the associated housing assembly, physically locating large conductive bars that extend the full length of the housing assembly is becoming increasingly difficult. Additionally, market demand for copper makes large ground bars expensive. By employing a short ground bar at the rear of the housing assembly, testing equipment can still effectively connect to, and provide the shunted path for, primary voltage and current; however, the requirement for a continuous connection throughout the draw-out range still must be met.\nThere is, therefore, a need for a housing assembly that provides for a grounded connection over the length of the draw-out path of travel but which does not include a long ground bar. There is a further need for such a housing assembly to accommodate existing electrical switching apparatus."}
{"text": "In U.S. patent application Ser. No. 509,770, filed July 1, 1983, now U.S. Pat. No. 4,500,657, for DENTAL RESTORATIVE COMPOSITIONS HAVING IMPROVED MECHANICAL PROPERTIES AND HYDROLYTIC STABILITY, and assigned to the same assignee as this Application, there is disclosed dental restorative compositions containing a monomer having at least two olefinically unsaturated groups, a polymerization initiator, and a filler. In the invention of said Application, both the monomer and the filler are hydrophobic (as defined therein), and the filler has a critical particle size and constitutes a critical proportion of the formulation. The dental restorative compositions of said application exhibit excellent properties, and are serious contenders to be indicated for use in Class II restorations (that is, restorations on the biting surfaces of molars), a distinction that has, thus far, eluded most, if not all, commercial resinous dental restoratives.\nThis invention is directed to an improvement in the art of resinous dental restoratives, and is particularly useful in the restoratives described in said Application."}
{"text": "This invention relates to the field of terminal clips and blocks, and connecting block assemblies thereof, used in telephone and data transmission interconnect assemblies. More particularly, this invention relates to a material reduced and transmission enhanced clip and block for clips of the type generally referred to in the art as 66 type clips. This invention also relates to a novel method of manufacturing such terminal clips.\nClip and block assemblies for 66 type clips are well known in the art and have been in use for many years. See, e.g., U.S. Pat. Nos. 3,112,147 and 3,518,618 to Pferd et al and to Swanson et al, respectively. These assemblies include a dielectric block or housing in which are located an ordered array of terminal clips. The clips are located in a series of openings in the block, and the top, or interconnect, end of each clip projects through the opening to extend above the upper surface of the central portion of the block. The bottom or base portion of each clip is retained within the body of the block by a combination of interaction with the block and a retainer plate which connects to the bottom of the block and retains the clips in position in the block.\nRecently, a requirement has arisen for a small 66 type clip that will be both (a) less expensive, because of reduced volume of material present in the clip, and (b) suitable for and capable of functioning in and for the same applications for which the previous 66 type clips have been used.\nOne attempt at a smaller and less expensive 66 type clip, and a block therefore, is disclosed in U.S. Pat. No. 5,127,845 to Ayer et al. According to the '845 patent, the thickness of the clip is reduced relative to a previous 66 type clip of U.S. Pat. No. 3,112,147, the upper portion of the clip is unchanged in configuration relative to that previous clip, and the base is reduced in size relative to that previous clip, so that the overall height of the clip is about three-fourths that of the previous clip and the clip requires about two thirds of the material of that previous clip. The '845 patent is the closest prior art to this invention known to the inventors of the subject matter of this application.\nAn important feature of the Ayer et al '845 patent is that the lower part of the clip has a downwardly divergent aperture (76) between the base and that portion of the clip which holds the wire, and that aperture is flanked by arm portions (78,80) of constant width. The aperture (76) is triangular in shape, and the opposed inner and outer surfaces of the arm portions (78 and 80) are parallel.\nAs with prior art patents, most notably the Swanson '618 patent, the terminal block assembly described by the Ayer et al '845 patent consists of a block body that positions clips at a constant spacing and a slidably engageable retaining plate designed to captivate the clips in said block body. An important limitation of these and other prior art patents is that the regular clip spacing causes uniform capacitive coupling between adjacent rows of clips such that, when used with twisted-pair wires, the capacitive coupling between \"tip\" and \"ring\" wires that constitute a pair is equal to that of adjacent conductors from different pairs. Since crosstalk performance is determined by the capacitive imbalance between pairs (i.e., the difference in capacitive coupling between each conductor of a pair and a conductor of another pair), the constant spacing between rows poses a limitation in terms of crosstalk performance between adjacent circuits.\nWhile it is appreciated that crosstalk performance may be improved by simply increasing the constant spatial separation between rows, such an approach has limitations and trade-offs in terms of diminished insertion loss and return loss performance caused by increased separation between tip and ring conductors of a pair, as well as limited compatibility with tools, adapters and other devices designed to plug directly into selected groups of clips on the connecting block assembly. Examples of such tools and plug-on devices are disclosed in U.S. Pat. Nos. 4,194,256, 4,585,290, 4,820,195, 4,883,430, 4,924,345 and 4,944,698, all of which are assigned to the assignee hereof.\nAll prior art blocks of this type are also limited by their height and the amount of material used to manufacture the connecting block body. The heights of the Pferd and Swanson blocks are determined by the height of the clip whereas the height of the Ayer et al block is determined by the combined height of the clip and that of the retainer ribs that support the clips. The disadvantage of these designs is that they are more limiting than other popular connecting block designs, such as the well known 110-series blocks disclosed in U.S. Pat. No. 4,964,812 (also assigned to the assignee hereof), in terms of the vertical clearance required for the block assembly when it is installed in telecommunications cabinets or other shallow telecommunications distribution facilities. The relatively greater height of the prior art block designs also requires more plastic material, which poses other limitations in terms of manufacturing cost and shipping weight."}
{"text": "The increasing complexity and compactness of present day electronic manufacturing provides increased packing density of various elements on semiconductor substrates and other electronic assemblies, such that a need exists for apparatus for holding and aligning preforms and especially preforms in the form of elongated elements, such as columns, or pins or the like, to facilitate appropriate mounting of these elements on similarly arrayed sites of an electronic component such as a substrate or package.\nIn the prior art, various devices for positioning discrete elements, such as solder balls, on electronic substrates have been provided. Typical of such arrangements, are devices which employ a vacuum for securing solder balls in a predetermined array in apertures of a principal face of a fixture for subsequent engagement of the solder balls with conductive pads, or other sites, on an electronic assembly.\nIn these prior art arrangements, the apparatus applies a vacuum to the fixture to provide a suction to secure a plurality of solder balls within apertures on a principal face of the fixture in an array that is complementary to select sites on an electronic substrate. Then, the fixture is appropriately positioned over the substrate with its arrayed solder balls essentially touching the selected sites of the electronic assembly, and subsequently upon termination of the vacuum, the solder balls are released from the fixture to their respective sites.\nWhile the above noted arrangements are well utilized in the construction of electronic devices, they still suffer from various disadvantageous. That is, as electronic assemblies continue to increase in complexity, additional needs such as more compact substrates and board designs are also required. With these requirements, users also simultaneously demand ever smaller and more portable devices that not only retain the capacity of larger devices but also provide increased capabilities. Hence, to meet these needs and demands, the packing density of electronic assemblies continues to rapidly increase. And as to be expected, with both increased capacity requirements and miniaturization, still ever greater burdens are placed upon the density and compactness of sites such as conductive pads on electronic substrates.\nIn order to satisfy this ever increasing density of sites on electronic substrates, an ever higher degree of precision is required in the placement of discrete elements, and especially elongated elements such as columns or pins on the substrate sites. In the prior art, solder balls are often rigidly held by the applied vacuum within conical or stepped apertures in the fixture's principal face. However, the discontinuities resulting from these conical or stepped apertures often result in offsetting, and in many cases, jamming of the preforms in their receiving apertures.\nMoreover, such discontinuities often trap debris which, in turn, will contribute to offsetting or jamming of the solder in their receiving apertures. These problems not only effect the precise positioning of the preforms in the fixture's array but, in addition to termination of the vacuum, often also require application of other mechanical forces to drive the tightly held solder balls from the fixture.\nIn turn, the difficulties of the above noted devices are further complicated where columns, including rods, pins, or other elongated preforms, rather than ball shaped preforms are utilized for mounting on electronic assemblies. For example, debris trapped in a conical, or stepped end of a column receiving hole can not only result in tipping and possible jamming of the column, but also can provide considerable variation in the relative position of their exterior ends."}
{"text": "The present invention relates to a semiconductor memory device having dynamic type RAM (DRAM) cells integrated and arranged, more specifically, to a semiconductor memory device capable of realizing a low power consumption, an enlarged operation margin and a high speed operation of a sense amplifier, a high reliability and a high current driving ability of a transistor by a transmission gate comprising a P-type MOS transistor.\nRecently, according to the increase of the demand for portable apparatus, smaller appliances mounted with a semiconductor device are demanded. Therefore, the mounted semiconductor device is required to be highly integrated and to consume low power so as to prolong the time to be used with a battery.\nAmong semiconductor memory devices, since a memory cell has a comparatively simple structure in a semiconductor memory device such as DRAM, a high integration has been achieved. According to the high integration, patterns and elements have been even minute to deteriorate the pressure resistance, and the like so that a lower voltage is needed for the power source voltage for driving according to the scaling rule so as to ensure the reliability and a lower power consumption.\nFor example, in a 16 M bit DRAM, the power source voltage is being shifted from the conventional Vcc=5V to Vcc=3.3V. Further, in a 256 M bit DRAM, the power source voltage to be supplied is considered to become Vcc=2.5V.\nFurthermore, as a minimum processing size, about 0.15 .mu.m is considered to be necessary. In a 1 G bit DRAM, a low power source voltage of about Vcc=1.8V is expected.\nIn these DRAMs, since a signal charge amount stored in a memory cell as data is minute, a highly sensitive bit line sense amplifier is needed for detecting, amplifying and reading it to the outside. Since the signal amount read out from the memory cell to be detected and amplified by the bit line sense amplifier is proportional to the power source voltage, when the power source voltage is scaled and lowered as mentioned above, the read signal amount decreases proportionally.\nIn general, as the bit line sense amplifier, a dynamic type differential amplifying circuit comprising an N-type sense amplifier where two sets of NMOS transistors (hereinafter NMOS is referred to as N-type) are connected by cross-coupling, and a P-type sense amplifier where two sets of PMOS transistors (hereinafter PMOS is referred to as P-type) are similarly connected by cross-coupling is used.\nAs shown in FIG. 28, as a precharge voltage of a bit line, which serves as the input terminal to the sense amplifier, the Vcc/2 precharge method, which is 1/2 of the power source voltage is used.\nTo explain the Vcc/2 precharge method simply, in a precharge cycle where the outside control signal RAS bar (hereinafter bar represents a reverse signal .sup.-) is at the \"H\" level, a bit line pair is precharged to Vcc/2. When the RAS bar signal shifts from the \"H\" to the \"L\" level to be in a state where the DRAM can read and write (active mode), the bit line precharge/equalize circuit is shut off so that the bit line precharged to be Vcc/2 becomes floating so that the word line WL selected by the external input address is driven, and data are read from the memory cell to the bit line.\nThe bit line sense amplifier is activated so that a minute signal read out to the bit line is detected and amplified.\nMore specifically, the \"L\" level side of the bit line pair is discharged to a ground voltage (Vss) by the above-mentioned N-type sense amplifier. On the other hand, the \"H\" level side of the bit line pair is charged to the power source voltage (Vcc) by the P-type sense amplifier. After completing the reading operation, the bit line is equalized and precharged to the Vcc/2 voltage again.\nAs a method for achieving a low power source voltage in a semiconductor memory device such as the above-mentioned DRAM, a method of reducing the bit line discharge current by having a smaller bit line voltage amplitude has been advocated conventionally.\nHowever, in the case with a small amplitude for the bit line voltage amplitude, since the voltage difference between a gate and a source of the transistor comprising the flip flop type sense amplifier becomes small, a problem of deterioration of the operation speed of the sense amplifier is involved. In a semiconductor memory device having a memory capacity of the G bit level, since the power source voltage Vcc becomes extremely small as 1.8V or less, a voltage between the gate and the source necessary for the sense amplifier transistor operation, that is, a voltage not less than the threshold voltage of the sense amplifier transistor cannot be obtained between the gate and the source of the sense amplifier transistor to disable the sense operation, and thus it is problematic.\nThat is, in the case of a DRAM using the above-mentioned Vcc/2 precharge method, due to the lower source voltage, in the sense amplifier operation, the voltage between the gate and the source (Vcc/2) to be applied to the transistor comprising the sense amplifier itself inevitably lowers. As a consequence, there is a risk that the sense operation is drastically delayed or the sense operation cannot be implemented.\nFor example, in the case of the source voltage Vcc=1.8V, only Vcc=0.9V is applied between the gate and the source of the sense amplifier transistor. Actually, the voltage is lowered by the resistance of the common source wiring of the sense amplifier transistor, in particular, in the initial sensing time, the value can be further smaller.\nOn the other hand, the absolute value .vertline.Vth.vertline. of the threshold voltage of the sense amplifier transistor (it is a positive voltage in an N-type sense amplifier, but a negative voltage in a P-type sense amplifier) needs to be at least 0.3V to 0.5V in order to guarantee the threshold value irregularity and the cut off characteristic.\nAt the time of the initial sensing, the substantial .vertline.Vth.vertline. is further raised in combination with the back gate effect of the sense amplifier transistor. Therefore, the voltage between the gate and the source of the sense amplifier transistor and the threshold voltage are extremely close at the time of the initial sensing so that the initial sensing is drastically delayed. And thus a serious problem is involved in terms of the high speed operation of the DRAM.\nAmong the transistors used in the DRAM, a highest voltage (boosting voltage) is applied to the gate electrode of the memory cell transistor. This is because the \"H\" level needs to be written in the memory cell. The gate voltage needed therefor is represented by VBLH+Vth'. Herein VBLH represents the \"H\" level voltage, and Vth' represents the threshold voltage of the memory cell transistor where a negative voltage is applied to the back gate (substrate) of the memory cell transistor and the source voltage is VBLH.\nIn order to maximize the charge to be accumulated in the memory cell capacitor in general, VBLH is equal to the source voltage in the chip (Vdd).\nIn general, since the gate oxide film of the same thickness is used for all the transistors in the chip for the production cost reduction, a transistor, of which gate electrode is not applied with the boosting voltage inevitably uses the thick gate oxide film the same as the memory cell transistor. Therefore, the transistors used in the DRAM have a problem in that their performance such as the current driving ability is low compared with a transistor of a logical semiconductor, and the like.\nIn order to improve the transistor performance by having a thinner gate oxide film, the boosting voltage can be lowered to reduce the voltage to be applied to the gate oxide film of the memory cell transistor. However, in that case, the source voltage in the chip inevitably becomes low so that since only 1/2 Vdd can be applied between the gate and the source of the transistor comprising the sense amplifier at most in the general 1/2 Vdd precharge method, a problem is involved in that the sense amplifier cannot function.\nIn general, an equalize circuit is provided in a bit line pair, a data line pair, a sense amplifier driving line, and the like. For example, a bit line equalize circuit conventionally comprises only N-type transistors as shown in FIG. 29.\nSince a large capacity is aimed at in a semiconductor memory these days, the redundancy technique for saving a defect memory cell generated in the production is in dispensable. However, even if it is saved, the defect cannot be eliminated.\nFor example, in the case a defect is generated in the short circuit between the bit line and the word line shown in FIG. 30, the memory cell is replaced (saved) by the redundancy technique. However, as mentioned above, since the defect remains as it is, at the time of operating the equalize circuit, through current flows in the route shown by the dot line in the figure, that is, VBL.fwdarw.PRCH.fwdarw.bit line.fwdarw.word line.fwdarw.row decoder.\nWhen the through current is large, or a large number of short circuit defects are generated between the bit line and the word line, although the memory cell can be saved with the redundancy technique, since the current consumption becomes larger than the standard range, the chip becomes \"defect\".\nIn order to restrain the through current, a method of inserting a current limiting circuit between the precharge circuit and the precharge voltage supply line as shown in FIG. 31 is proposed. It is preferable that the current limiting circuit comprises a P-type transistor. At the time, problems are involved in that (1) the sense amplifying area is enlarged due to the necessity of the well separation with respect to the equalize/precharge circuit portion, and (2) a wiring is required for connecting the precharge circuit and the current limiting circuit."}
{"text": "The present invention relates generally to the field of luggage and equipment for carrying luggage and more particularly, to an improved shoulder pad.\nThe prior art related to luggage and hand-carried equipment which includes various types of shoulder straps and shoulder pads which attempt to decrease the discomfort associated with carrying luggage or equipment by using a shoulder pad to distribute the load exerted by the weight of the equipment. Typically, conventional shoulder pads consist of a generally flat rectangular pad which is made of various types of rubber and a means for attaching the pad to a shoulder strap.\nConventional shoulder pads typically do not accomplish the goals of reducing the discomfort which results from the use of a shoulder strap. Conventional shoulder pads flex easily and do not adequately spread the load exerted on the shoulder of a user by the shoulder strap and thus do not adequately reduce the load per unit of area of contact between the shoulder pad and the user\"\"s shoulder. As a result, the conventional shoulder pad digs into the user\"\"s shoulder and severely limits both the load which can be carried and the length of time that the load can be carried before the degree of user discomfort becomes unacceptable.\nIn addition, a shoulder strap, when used with a conventional shoulder pad, tends to slide off the shoulder of a user. As a result, a user is forced to interrupt walking or other activity and continually readjust the position of the shoulder strap. This need for continuing periodic adjustment creates a degree of discomfort which further limits the utility of conventional shoulder pads.\nIt is an object of the present invention to provide a shoulder pad which effectively prevents the shoulder strap from digging into the shoulder of a user.\nAnother object of the present invention is to provide a shoulder pad which can be used with relatively heavy loads for projected periods of time.\nAnother object of the present invention is to provide a shoulder pad which prevents the shoulder strap from slipping off the shoulder of a user.\nAnother object of the present invention is to provide a high degree of comfort during use.\nAnother object of the present invention is to provide a shoulder pad which can be quickly and easily installed on a shoulder strap.\nYet another object of the present invention is to provide a shoulder pad which comprises a relatively small number of component parts which can be manufactured economically in volume resulting in a relatively low unit cost.\nThe foregoing and other objects and advantages of the invention will appear more clearly hereinafter.\nIn accordance with the present invention, a shoulder pad includes a relatively rigid body member and a relatively soft pad member. The pad member is mounted on a curved lower surface of the body member. The upper surface of the body member includes a slot which accepts a luggage strap and which is angled forwardly toward the lower surface. The curved lower surface accommodates the curvature of the user\"\"s shoulder while the pad member conforms to the individual user\"\"s shoulder anatomy. The upper and lower surfaces of the body member are formed in an angular relationship so that during use, the slot is always directed toward the user\"\"s shoulder so that the load exerted by the shoulder strap keeps the shoulder pad in place on the user\"\"s shoulder.\nIn an alternative embodiment of the invention, the body member is formed of a relatively flexible plastic material and includes a central portion which has an upper surface disposed at an angle relative to a lower surface. A pair of clip members facilitate attachment of the body member to a shoulder strap. During use, the shoulder strap extends over the upper surface of the central portion and under the clips with the angular central portion forcing the shoulder strap into the closed portion of the clips. The flexible body member bends to conform to the user\"\"s shoulder.\nA relatively soft pad member is mounted on the lower surface of the body member. The pad member increases the comfort of the body member on the user\"\"s shoulder. A central portion of the pad member projects into an aperture which is formed in the central portion of the body member thereby aiding in locking the pad member and the body member together."}
{"text": "There is an increasing desire to develop sustainable fuel options that are not reliant of fossil fuels. One option being evaluated is the production of biofuels obtained from algae. The commercialization of such biofuels will likely require cultivating algae in large ponds. The ponds of algae would then be harvested and the lipids within the algae extracted and converted into biofuel. Existing methods of harvesting the algae have not proven commercially viable for numerous reasons, which include, but are not limited to: (i) too much water remains in the extracted algae for efficient lipid extraction, (ii) the algae cell walls are ruptured prematurely causing release of the valuable lipids before the extraction and conversion processes, and (iii) the algae cell walls are ruptured preventing the algae from being reused. Thus, there is a need for improvements related to the harvesting and processing of algae for the production of biofuels."}
{"text": "I. Technical Field\nThe present invention relates to twin-skeg ships having a pair of left and right skegs provided on the bottom of the stern so as to support propeller shafts.\nII. Description of the Related Art\nA twin-skeg ship having a pair of left and right propeller shafts supported by skegs integrated with the hull so as to form a tunnel-like bottom recess between the left and right skegs is known in the related art (see, for example, Japanese Unexamined Patent Application, Publication No. HEI-8-133172).\nOn the other hand, a reaction fin has been devised as a propulsion-performance improving unit installed in a ship to improve the propulsion performance of the ship (unit for improving the speed of the ship with the same horsepower) (see, for example, Japanese Unexamined Patent Application, Publication No. HEI-5-185986)."}
{"text": "Large and/or complicated computer networks may contain multiple routing devices and multiple bridging devices within the topmost level of the network hierarchy. Such a network can be considered to consist of several top-level routing realms. A routing realm is, within the context of this document, a term referring to the network space behind a routing device, or the network space behind a group of interconnected bridging devices.\nFor example, a network containing two top-level routers R1 and R2, and four top-level bridging devices B1, B2, B3, and B4, can be considered to have three top-level routing realms: R1, R2, and (B1, B2, B3, B4).\nIn the case of a network which is provisioned by a single device at the top-level of the network, the network's routing realms, by virtue of fundamental networking principles, must utilize a distinct collection of Internet Protocol (IP) addresses, relative to the other routing realms within that network.\nA network consisting of several top-level routing realms, where at least one of those routing realms is defined by a routing device, presents challenges to a top-level network provisioning device, because the provisioning device must successfully process routed and bridged traffic simultaneously. Such routed traffic is ISO Open Systems Interconnection Layer-3 traffic and such bridged traffic is ISO OSI Layer-2 traffic.\nWhen an end-user device moves from one routing realm to another, since each routing realm must utilize a relatively unique collection of IP addresses, the IP address of the end-user device must change as the end-user device enters into a new routing realm. Thus, the existing provisioning device cannot handle such user who moves from a routing realm to another.\nIt is therefore desirable to provide a network provisioning device that can recognize that an end-user device has moved from a routing realm to another in a network comprising multiple routing realms, and handle this situation automatically."}
{"text": "The present invention relates to thermal ink jet printers, and more particularly, to printheads incorporating a plurality of heater resistors which are selectively addressed to heat and expel ink from adjacent ink channels.\nThermal ink jet printing utilizes printheads which use thermal energy to produce vapor bubbles in ink-filled channels to expel ink droplets. A thermal energy generator or heater element, usually a resistor, is located at a predetermined distance from a nozzle of each one of the channels. The resistors are individually addressed with an electrical pulse to generate heat which is transferred from the resistor to the ink.\nThe transferred heat causes the ink to be super heated, i.e., far above the ink's normal boiling point. For example, a water based ink reaches a critical temperature of 280.degree. C. for bubble nucleation. The nucleated bubble or water vapor thermally isolates the ink from the heater element to prevent further transfer of heat from the resistor to the ink. Further, the nucleating bubble expands until all of the heat stored in the ink in excess of the normal boiling point diffuses away or is used to convert liquid to vapor which, of course, removes heat due to heat of vaporization. During the expansion of the vapor bubble, the ink bulges from the nozzle and is contained by the surface tension of the ink as a meniscus.\nWhen the excess heat is removed from the ink, the vapor bubble collapses on the resistor, because the heat generating current is no longer applied to the resistor. As the bubble begins to collapse, the ink still in the channel between the nozzle and bubble starts to move towards the collapsing bubble, causing a volumetric contraction of the ink at the nozzle and resulting in the separating of the bulging ink as an ink droplet. The acceleration of the ink out of the nozzle while the bubble is growing provides the momentum and velocity to expel the ink droplet towards a recording medium, such as paper, in a substantially straight-line direction. The entire bubble expansion and collapse cycle takes about 20 microseconds (.mu.s). The channel can be refired after 100 to 500 .mu.s minimum dwell time to enable the channel to be refilled and to enable the dynamic refilling factors to be somewhat dampened.\nTo eject an ink droplet, each heater element must become hot enough to cause the ink to reach a bubble nucleation temperature of at least 280.degree. C. for water based ink. In order for the heater element to generate the thermal energy to cause bubble nucleation, an operating voltage is applied to a resistor of the heater element. Typically, the operating voltage is proportional to the resistance of the resistor, i.e., the higher the resistance, the higher the operating voltage.\nConventionally, polysilicon is used to form the resistors of the heater elements. The resistance value of the resistors is chosen based on the actual required power (Power=V.times.I=I.sup.2 .times.R=V.sup.2 /R) for ejection of the ink droplet through bubble nucleation. Once the required power and voltage has been chosen, the resistance value is determined. The fabrication of the determined resistance is controlled by the sheet resistance (ohms/square; .OMEGA./.quadrature.) of the polysilicon and the size of the resistor. The size of the resistor can be tightly controlled by photolithography and polysilicon etch techniques. The sheet resistance of the polysilicon is primarily controlled by impurity doping, preferably by ion implantation, and annealing of the doped polysilicon.\nAs drop ejection becomes more efficient, the heater element size required decreases and the resistance required increases. The more resistive the heaters, the lower the doping level, and the more difficult it is to accurately control the resistance. When resistance is too high, threshold voltage is too high and drops are not ejected. When resistance is too low, threshold voltage is too low, and the extra voltage applied to the printhead results in excess heat, which causes ink to bake on the heater. This results in an extra insulting layer on the heater, decreasing energy transferred to the ink, creating smaller drops and eventually completely failing to eject drops. In addition, \"over voltage\" operation results in decreased heater lifetime. Conventional techniques use doping to introduce impurities in the polysilicon heater structure. One method disclosed in U.S. Pat. No. 5,639,386, whose contents are incorporated by reference, teaches methods of fabricating polysilicon resistors with improved threshold uniformity which includes steps for forming resistors with two doped regions, a lightly doped center region and more heavily doped end region. The heater element is formed on a polysilicon layer 8 (see FIG. 3 of the patent). N-type dopants, typically phosphorus, are ion implanted into the polysilicon layer to form a lightly doped heater center region. A photoresist layer 8 is patterned to protect the heater center region followed by heavier doping at the ends. This results in lightly doped center regions 8A and heavily doped end regions 8B. Regions 8B are heavily doped to make good electrical contacts to address leads and for low resistance to form transistor gates. The photoresist layer is removed, and the polysilicon is plasma etched to form the final resistor structure."}
{"text": "The present invention concerns a system for resampling, and in particular, to a digital resampling system which resamples signals using a quantized phase.\nSample re-conversion systems are well known for converting information bearing signals from one format to another format. Such systems are commonly used for data telemetry, audio signal processing, video signal processing, and video signal standards conversion.\nAn exemplary video signal processing system is described in U.S. Pat. No. 4,774,581 entitled TELEVISION PICTURE ZOOM SYSTEM, issued to Shiratsuchi, which is hereby incorporated by reference for its teachings on digital resampling techniques.\nAnother exemplary video system standards conversions system is described in U.S. Pat. No. 5,057,911 entitled, SYSTEM AND METHOD FOR CONVERSION OF DIGITAL VIDEO SIGNALS, which is hereby incorporated by reference for its teachings on digital resampling techniques.\nThe system described in the SYSTEM AND METHOD FOR CONVERSION OF DIGITAL VIDEO SIGNALS patent converts video signals having a normal sample rate of 13.5 MHz into a digital signal having a sample rate of 14.31818 MHz, (hereinafter 14.3 MHz) which is compatible with the National Television Standards Committee (NTSC) standard. The ratio of these two sample frequencies is exactly equal to the ratio of 33 to 35.\nAccordingly, the standards conversions system described in the referenced patent generates 35 output samples for every 33 input samples it receives. This is accomplished using 35 interpolation filters which generate interpolated samples at 35 respective positions between any two of the input samples. The filter, X(i+1), which is selected to generate the i+1.sup.th output sample is defined by equation (1). EQU X(i+1)=(X(i)+33) MOD 35 (1)\nIn equation (1) 33 is the number of input samples, 35 is the number of output samples, MOD is an abbreviation of a modulo operation, X is the interpolation phase location for the current sample in the sequence, and i is the time index. The interpolation phase location defines a particular frequency response characteristic to be used to filter the input samples in order to generate the particular resampled output sample.\nThis technique of sample rate conversion works well when the input and output sample rates can be represented as a ratio of two relatively small numbers. When, however, the ratio of these rates cannot be represented to a high degree of precision by a manageable ratio, this conversion process may be inaccurate or unduly expensive to implement.\nFor example, consider a resampling circuit which converts from the CCIR (601) standard having a nominal sample rate of 13.5 MHz to the Phase Alternate Line (PAL) standard which has a sample rate of 17.734475 (hereinafter 17.7) MHz. The smallest exact ratio between these two sample frequencies is 540,000 over 709,379. Thus, if the scheme described above were used to resample the CCIR (601) signal in a PAL digital signal, 709,379 digital filters would be required. At the current state of the art, it is not practical to implement a video standards conversion system having 709,730 filters.\nAn alternative sampling scheme may be used where the resampling function is approximated by limiting the number of possible interpolation points. As a result, the actual interpolated point is estimated by choosing the filter that is closest to the desired interpolated point. This is known as sample phase quantization. This technique introduces a phase shift error in the signal. However, with proper quantization, this error may be acceptable.\nA system employing sample phase quantization would include a fixed number of filters or a filter with a fixed number of frequency characteristics which number is designated as NPHASES. Next, the number of filters in the ideal system is quantized into the number of filters, NPHASES, for the system. This is accomplished using equations (2) and (3) below. EQU X.sub.i+1 =(X.sub.i +A) MOD B (2) EQU Q=X.sub.i+1 * NPHASES/B (3)\nIn equations (2) and (3), B is the number of samples in the output line of the filter, A is the number of samples in the input line of the filter, and Q is the quantized phase. The quantized phase is used to determine which filter in the system is to be selected to produce the new sample.\nAlthough equations (2) and (3) may be used to quantize the filters and phases, the equations require division and multiplication operations which are computationally expensive."}
{"text": "The invention had the object of finding novel compounds having valuable properties, in particular those which can be used for the preparation of medicaments.\nThe present invention relates to compounds and to the use of compounds in which the inhibition, regulation and/or modulation of signal transduction by kinases, in particular tyrosine kinases and/or serine/threonine kinases, plays a role, furthermore to pharmaceutical compositions which comprise these compounds, and to the use of the compounds for the treatment of kinase-induced diseases.\nIn particular, the present invention relates to compounds and to the use of compounds in which the inhibition, regulation and/or modulation of signal transduction by Met kinase plays a role.\nOne of the principal mechanisms by which cellular regulation is effected is through the transduction of extracellular signals across the membrane that in turn modulate biochemical pathways within the cell. Protein phosphorylation represents one course by which intracellular signals are propagated from molecule to molecule resulting finally in a cellular response. These signal transduction cascades are highly regulated and often overlap, as is evident from the existence of many protein kinases as well as phosphatases. Phosphorylation of proteins occurs predominantly at serine, threonine or tyrosine residues, and protein kinases have therefore been classified by their specificity of phosphorylation site, i.e. serine/threonine kinases and tyrosine kinases. Since phosphorylation is such a ubiquitous process within cells and since cellular phenotypes are largely influenced by the activity of these pathways, it is currently believed that a number of disease states and/or diseases are attributable to either aberrant activation or functional mutations in the molecular components of kinase cascades. Consequently, considerable attention has been devoted to the characterisation of these proteins and compounds that are able to modulate their activity (for a review see: Weinstein-Oppenheimer et al. Pharma. &. Therap., 2000, 88, 229-279).\nThe role of the receptor tyrosine kinase Met in human oncogenesis and the possibility of inhibition of HGF (hepatocyte growth factor) dependent Met activation are described by S. Berthou et al. in Oncogene, Vol. 23, No. 31, pages 5387-5393 (2004). The inhibitor SU11274 described therein, a pyrrole-indoline compound, is potentially suitable for combating cancer. Another Met kinase inhibitor for cancer therapy is described by J. G. Christensen et al. in Cancer Res. 2003, 63(21), 7345-55.\nA further tyrosine kinase inhibitor for combating cancer is reported by H. Hov et al. in Clinical Cancer Research Vol. 10, 6686-6694 (2004). The compound PHA-665752, an indole derivative, is directed against the HGF receptor c-Met. It is furthermore reported therein that HGF and Met make a considerable contribution to the malignant process of various forms of cancer, such as, for example, multiple myeloma.\nThe synthesis of small compounds which specifically inhibit, regulate and/or modulate signal transduction by tyrosine kinases and/or serine/threonine kinases, in particular Met kinase, is therefore desirable and an aim of the present invention.\nIt has been found that the compounds according to the invention and salts thereof have very valuable pharmacological properties while being well tolerated.\nThe present invention specifically relates to compounds of the formula I which inhibit, regulate and/or modulate signal transduction by Met kinase, to compositions which comprise these compounds, and to processes for the use thereof for the treatment of Met kinase-induced diseases and complaints, such as angiogenesis, cancer, tumour formation, growth and propagation, arteriosclerosis, ocular diseases, such as age-induced macular degeneration, choroidal neovascularisation and diabetic retinopathy, inflammatory diseases, arthritis, thrombosis, fibrosis, glomerulonephritis, neurodegeneration, psoriasis, restenosis, wound healing, transplant rejection, metabolic diseases and diseases of the immune system, also autoimmune diseases, cirrhosis, diabetes and diseases of the blood vessels, also instability and permeability and the like in mammals.\nSolid tumours, in particular fast-growing tumours, can be treated with Met kinase inhibitors. These solid tumours include monocytic leukaemia, brain, urogenital, lymphatic system, stomach, laryngeal and lung carcinoma, including lung adenocarcinoma and small-cell lung carcinoma.\nThe present invention is directed to processes for the regulation, modulation or inhibition of Met kinase for the prevention and/or treatment of diseases in connection with unregulated or disturbed Met kinase activity. In particular, the compounds of the formula I can also be employed in the treatment of certain forms of cancer. The compounds of the formula I can furthermore be used to provide additive or synergistic effects in certain existing cancer chemotherapies, and/or can be used to restore the efficacy of certain existing cancer chemotherapies and radiotherapies.\nThe compounds of the formula I can furthermore be used for the isolation and investigation of the activity or expression of Met kinase. In addition, they are particularly suitable for use in diagnostic methods for diseases in connection with unregulated or disturbed Met kinase activity.\nIt can be shown that the compounds according to the invention have an antiproliferative action in vivo in a xenotransplant tumour model. The compounds according to the invention are administered to a patient having a hyperproliferative disease, for example to inhibit tumour growth, to reduce inflammation associated with a lymphoproliferative disease, to inhibit trans-plant rejection or neurological damage due to tissue repair, etc. The present compounds are suitable for prophylactic or therapeutic purposes. As used herein, the term “treatment” is used to refer to both prevention of diseases and treatment of pre-existing conditions. The prevention of proliferation is achieved by administration of the compounds according to the invention prior to the development of overt disease, for example to prevent the growth of tumours, prevent metastatic growth, diminish restenosis associated with cardiovascular surgery, etc. Alternatively, the compounds are used for the treatment of ongoing diseases by stabilising or improving the clinical symptoms of the patient.\nThe host or patient can belong to any mammalian species, for example a primate species, particularly humans; rodents, including mice, rats and hamsters; rabbits; horses, cows, dogs, cats, etc. Animal models are of interest for experimental investigations, providing a model for treatment of human disease.\nThe susceptibility of a particular cell to treatment with the compounds according to the invention can be determined by in vitro tests. Typically, a culture of the cell is combined with a compound according to the invention at various concentrations for a period of time which is sufficient to allow the active agents to induce cell death or to inhibit migration, usually between about one hour and one week. In vitro testing can be carried out using cultivated cells from a biopsy sample. The viable cells remaining after the treatment are then counted.\nThe dose varies depending on the specific compound used, the specific disease, the patient status, etc. A therapeutic dose is typically sufficient considerably to reduce the undesired cell population in the target tissue while the viability of the patient is maintained. The treatment is generally continued until a considerable reduction has occurred, for example an at least about 50% reduction in the cell burden, and may be continued until essentially no more undesired cells are detected in the body.\nFor identification of a signal transduction pathway and for detection of interactions between various signal transduction pathways, various scientists have developed suitable models or model systems, for example cell culture models (for example Khwaja et al., EMBO, 1997, 16, 2783-93) and models of transgenic animals (for example White et al., Oncogene, 2001, 20, 7064-7072). For the determination of certain stages in the signal transduction cascade, interacting compounds can be utilised in order to modulate the signal (for example Stephens et al., Biochemical J., 2000, 351, 95-105). The compounds according to the invention can also be used as reagents for testing kinase-dependent signal transduction pathways in animals and/or cell culture models or in the clinical diseases mentioned in this application.\nMeasurement of the kinase activity is a technique which is well known to the person skilled in the art. Generic test systems for the determination of the kinase activity using substrates, for example histone (for example Alessi et al., FEBS Lett. 1996, 399, 3, pages 333-338) or the basic myelin protein, are described in the literature (for example Campos-González, R. and Glenney, Jr., J. R. 1992, J. Biol. Chem. 267, page 14535).\nFor the identification of kinase inhibitors, various assay systems are available. In scintillation proximity assay (Sorg et al., J. of Biomolecular Screening, 2002, 7, 11-19) and flashplate assay, the radioactive phosphorylation of a protein or peptide as substrate with γATP is measured. In the presence of an inhibitory compound, a decreased radioactive signal, or none at all, is detectable. Furthermore, homogeneous time-resolved fluorescence resonance energy transfer (HTR-FRET) and fluorescence polarisation (FP) technologies are suitable as assay methods (Sills et al., J. of Biomolecular Screening, 2002, 191-214).\nOther non-radioactive ELISA assay methods use specific phospho-antibodies (phospho-ABs). The phospho-AB binds only the phosphorylated substrate. This binding can be detected by chemiluminescence using a second peroxidase-conjugated anti-sheep antibody (Ross et al., 2002, Biochem. J.).\nThere are many diseases associated with deregulation of cellular proliferation and cell death (apoptosis). The conditions of interest include, but are not limited to, the following. The compounds according to the invention are suitable for the treatment of various conditions where there is proliferation and/or migration of smooth muscle cells and/or inflammatory cells into the intimal layer of a vessel, resulting in restricted blood flow through that vessel, for example in the case of neointimal occlusive lesions. Occlusive graft vascular diseases of interest include atherosclerosis, coronary vascular disease after grafting, vein graft stenosis, peri-anastomatic prosthetic restenosis, restenosis after angioplasty or stent placement, and the like."}
{"text": "It is difficult for a computer-implemented process to identify the part of a song that a listener would find interesting. For example, a computer process may receive a waveform of a song. However, the computer process may not be able to identify which part of the song a listener would find interesting or memorable."}
{"text": "1. Field of the Invention\nThe present invention relates generally to non-volatile magnetic memory and particularly to non-uniform switching of non-volatile magnetic based memory.\n2. Description of the Prior Art\nComputers conventionally use rotating magnetic media, such as hard disk drives (HDDs), for data storage. Though widely used and commonly accepted, such media suffer from a variety of deficiencies, such as access latency, higher power dissipation, large physical size and inability to withstand any physical shock. Thus, there is a need for a new type of storage device devoid of such drawbacks.\nOther dominant storage devices are dynamic random access memory (DRAM) and static RAM (SRAM) which are volatile and very costly but have fast random read/write access time. Solid state storage, such as solid-state-nonvolatile-memory (SSNVM) devices having memory structures made of NOR/NAND-based Flash memory, providing fast access time, increased input/output (IOP) speed, decreased power dissipation and physical size and increased reliability but at a higher cost which tends to be generally multiple times higher than hard disk drives (HDDs).\nAlthough NAND-based flash memory is more costly than HDD's, it has replaced magnetic hard drives in many applications such as digital cameras, MP3-players, cell phones, and hand held multimedia devices due, at least in part, to its characteristic of being able to retain data even when power is disconnected. However, as memory dimension requirements are dictating decreased sizes, scalability is becoming an issue because the designs of NAND-based Flash memory and DRAM memory are becoming difficult to scale with smaller dimensions. For example, NAND-based flash memory has issues related to capacitive coupling, few electrons/bit, poor error-rate performance and reduced reliability due to decreased read-write endurance. Read-write endurance refers to the number of reading, writing and erase cycles before the memory starts to degrade in performance due primarily to the high voltages required in the program, erase cycles.\nIt is believed that NAND flash, especially multi-bit designs thereof, would be extremely difficult to scale below 45 nanometers. Likewise, DRAM has issues related to scaling of the trench capacitors leading to very complex designs which are becoming increasingly difficult to manufacture, leading to higher cost.\nCurrently, applications commonly employ combinations of EEPROM/NOR, NAND, HDD, and DRAM as a part of the memory in a system design. Design of different memory technology in a product adds to design complexity, time to market and increased costs. For example, in hand-held multi-media applications incorporating various memory technologies, such as NAND Flash, DRAM and EEPROM/NOR flash memory, complexity of design is increased as are manufacturing costs and time to market. Another disadvantage is the increase in size of a device that incorporates all of these types of memories therein.\nThere has been an extensive effort in development of alternative technologies such as Ovanic Ram (or phase-change memory), Ferromagnetic Ram (FeRAM), Magnetic Ram (MRAM), Nanochip, and others to replace memories used in current designs such as DRAM, SRAM, EEPROM/NOR flash, NAND flash and HDD in one form or another. Although these various memory/storage technologies have created many challenges, there have been advances made in this field in recent years. MRAM seems to lead the way in terms of its progress in the past few years to replace all types of memories in the system as a universal memory solution.\nA prior art field-switching MRAM structure used in conventional MRAMs is depicted in FIG. 1. In FIG. 1, an MRAM cell 10 is shown to include a bit line 12 and a transistor 14 and formed there between is the memory element 32 and a number of metal lines 20, 22, 24, 26 formed for ease of manufacturing. The word line (WL) 16 is shown formed on top of the gate of the transistor 14 and the digit line (DL) is shown formed on the bottom of the memory element 32. The memory element consists of three key layers namely the fixed layer, the barrier tunneling layer and the free-layer. In operation, the digit line (DL) 14 is used to change the magnetic orientation of the free-layer of the memory element and thereby creating “parallel” (low resistance) and “anti-parallel” (high resistance) states which become the “0” and “1”.\nIn FIG. 1, M1s 20 and 28, V1 22, M2 24 and V2 26 are examples of how these layers would be processed in order to connect the memory element (or cell) to the transistor 14. The M1 28, on the transistor 14's “drain” side, is connected to a common ground and is deposited, at substantially the same time, as the M1 20 on the “source” side of the transistor 14. The bit line 12 is typically connected through the memory element 32 to the source 30 of the transistor 14. The word-line (WL) 16 is typically connected to the control gate of the transistor 14 for selecting the specific transistor.\nThe problem with the MRAM cell 10 of FIG. 1 are two-fold namely, the large cell-size and therefore high cost and high-power. The switching of the memory element 32 is made through the magnetic-field generated from the digit-line 18 which limits how closely the neighboring cells can be and thereby leading to larger cell size and thereby higher cost.\nOne of the problems with prior art memory structures is that the current and power requirements are too high to make a functional memory device. This also poses a key concern regarding the reliability of such device due to likely dielectric break-down of the tunneling barrier layer and thereby making it non-functional.\nFIG. 2 shows relevant layers of the memory elements 32 of prior art (MRAMs), such as the cells 10 (FIG. 1) and memory element 32. In FIG. 2, a free layer 60 is shown on top of which is shown formed a tunnel layer 62 on top of which is shown formed a fixed layer 64. The free layer's 60 magnetic moment can change, whereas, the moment of the fixed layer 64, below a known temperature, remains fixed. The tunneling layer 62 is commonly referred to as the “barrier layer”. In some prior art structures, the free layer 60 is made of several free layers. Application of current to the structure of FIG. 2 causes switching between anti-parallel (AP) and parallel (P) states, which, in turn represent two logical states for storing information in memory made from the structure of FIG. 2. The relationship between resistances of the two states depends on the tunneling-magneto-resistance (TMR) is defined as:TMR=(Rh−Rl)/Rl  Eq. (1)\nWherein Rh is resistance at a high state and Rl is resistance at a low state.\nLow capacity MRAM memory based on a design relying on magnetic-field to switch the memory elements is another known memory. It has been shown that current can also be used to switch the memory elements. The challenge has been that the switching current is too high to allow the making of a functional device for memory applications due to the memory's high power consumption. Several recent publications, such as those cited below as references 5 and 6(5,6) have shown that the switching current can be reduced by having the memory element pinned by two anti-ferromagnetic (AF)-couple layers resulting in spin oscillations or “pumping” and thereby reducing the switching current.\nWhat is needed is storage memory based on magnetic memory for storage of digital information and having reduced switching current in the magnetic memory thereby decreasing power consumption and reduced cell size thereby reducing costs associated with manufacturing the memory."}
{"text": "In order to comply with ever more stringent limitations on pollutant emissions, a plurality of actions are taken in present-day engines in order to reduce, in particular, emissions of particulates and of nitrogen oxides.\nOne essential action is exhaust gas recirculation (EGR), which represents a present-day arrangement for preventing nitrogen oxide emissions. Exhaust gas recirculation lowers the oxygen content in the cylinder, and a decrease in the temperature in the combustion chamber consequently occurs. A rise in particulates with increasing exhaust gas recirculation is problematic. The principal reason for higher particulate emissions is that the oxygen also necessary for carbon oxidation is limited. The reduction in oxygen content caused by exhaust gas recirculation thus always has the effect of decreasing nitrogen oxide emissions, and elevating particulate emissions. This creates a conflict of objectives between particulate and nitrogen-oxide emissions, especially in diesel engines.\nAs a result of existing regulatory stipulations for the exhaust test cycle, the requirements for passenger cars regarding reduction of pollutant emissions in dynamic operation have been undemanding. In the commercial vehicle sector, dynamic operation has been entirely excluded by a steady-state test.\nThe development of exhaust gas recirculation control systems has therefore been directed essentially toward reducing emissions under steady-state conditions. The EGR control structure has therefore been based on values for engine speed and engine load that were identified under steady-state conditions and stored in characteristics diagrams. Control of air masses or EGR rates is known from known EGR control systems. Future regulatory requirements for commercial-vehicle and passenger-vehicle engines provide for a substantially larger dynamic component. In the future, certification will focus on emissions occurring in actual driving operation, e.g. in dynamic operation (“real driving emissions” or RDE), and on fuel consumption.\nTaking dynamic processes into account requires, in particular, taking into account abrupt load changes and rapid load increases, such as those that often occur in real driving operation and in future test cycles. In a diesel engine, abrupt load changes and rapid load increases result, because of the inertia of the air system, in a delay in boost pressure buildup. The reasons for this inertia include the moment of inertia of the turbocharger and the dead volume between the compressor and the engine intake valves. The injection system that implements the driver's load demand has a considerably shorter reaction time than the engine's air system.\nBecause cylinder filling in a diesel engine is determined substantially by the slow-reacting boost pressure, the target values of the air system—based on injection volume and engine speed—are mismatched with the dynamic state of the engine. The steady-state target value of an air mass control system thus results, in the context of a slow-reacting boost pressure buildup and thus decreased cylinder filling, in a sharp reduction in the EGR rate, thereby producing dynamic nitrogen oxide spikes. With an EGR rate control system, the result is a smaller air mass and thus elevated particulate emissions. In systems or operating states dominated by high-pressure EGR, the resulting additional negative effect on boost pressure buildup potentially limits the maximum possible injection volume, with a consequent delay in reaching the desired target torque.\nPatent document DE 100 10 978 A1 discusses a method for regulating the boost pressure of an internal combustion engine having an exhaust gas-driven turbocharger whose turbine, disposed in the exhaust duct of the internal combustion engine, has a modifiable geometry, regulation of the boost pressure being accomplished by way of an adjustment of the turbine geometry. In order to ensure boost regulation such that in the context of a load change, the boost pressure tracks as quickly as possible the change in the desired target boost pressure value, this method and this apparatus provide for determining a control variable for the turbine geometry from the system as a function of the exhaust backpressure existing in the exhaust duct in front of the turbine. This is because the exhaust backpressure reacts considerably more quickly than the boost pressure to a change in the behavior of the controlled system, e.g. to a load alternation, change in engine speed, change in exhaust gas recirculation, or to malfunctions, e.g. in the actuation system.\nPatent document DE 41 07 693 A1 discusses a system for regulating and controlling a turbocharger, in which a predefined target value of the boost pressure is compared with an actual value. As a function of the comparison, a control system generates an actuating signal for applying control to an actuating mechanism. The target value depends at least on an engine speed signal and on a load signal. An arrangement is provided which elevates the boost pressure upon a specific change in the load signal.\nThe methods from the existing art for identifying target values for the control function are not sufficient to react, in the control loops for exhaust gas recirculation, turbocharging, and fuel injection, to the state of affairs described above in transient driving situations. In methods known from the existing art, transient driving situations of this kind are implemented, for example, by monitoring accelerator pedal gradients or injection volume gradients. The problem here is that these variables are not relevant in terms of preventing the generation of pollutant emissions.\nStrictly in principle, it is also possible to use electric motors (e-machines) or electrical auxiliary compressors (e-boosters) to accelerate the buildup of boost pressure or torque. Here as well, a transient driving situation must be optimally detected."}
{"text": "Nonflammable injection-molding resins composed of a low-molecular-weight styrene resin and a low-polymerization-degree vinyl chloride resin are superior materials in view of mechanical properties and cost. Recently, these resins have been used in increased amounts as materials for housings of office automation (OA) apparatuses and the like. However, the requirements of molded articles are becoming more severe each year in OA apparatus applications. In particular, miniaturization of OA apparatuses, because it brings about an increase in internal heat generation in the apparatuses, necessitates a decrease in the thickness of the housing materials, and requires the use of highly heat-resistant resins. A general means for raising heat-resistance is to use a styrene resin which contains .alpha.-methylstyrene as a third component, or to use a post-chlorinated vinyl chloride resin.\nThese resins, however, are disadvantageously liable to cause jellowing and flashing of the molded articles and are not readily moldable. The disadvantages of the molded articles are considered to relate to the extremely high melt viscosity of the resins in injection molding, particularly in the case of a resin which contains a post-chlorinated vinyl chloride type resin. The use of such a resin results in large heat generation by the shearing action of metal parts such as a screw, a cylinder, a mold, etc. and causes heat decomposition of the resin. For improving the molding of such resins, methods such as the addition of a large amount of a stabilizer or a lubricant, and the addition of a moldability-improving agent comprising an alkyl acrylate type resin are generally used. The former method is not preferable, because the cost of the starting materials increases and, further, the impact strength and the heat distortion temperature are decreased adversely affecting the quality of the molded articles. The latter method, namely the addition of an alkyl acrylate type resin as a moldability-improving agent, is remarkably useful since the increase in the cost of the material is not so large and lowering of the heat distortion temperature is not remarkable. This method, however, causes a decrease of the impact strength and peeling of the molded articles. In particular, with an increase in the amount of the additive for increasing the moldability-improving effect, lowering of the impact strength and deterioration of the non-peeling property of the molded articles become remarkable. These are serious problems.\nThese problems are assumed to be caused by the alkyl acrylate type resin which is used as a moldability-improving agent and which is designed to improve moldability particularly of extrusion-molding vinyl chloride resins having an ordinary molecular weight. The inventors of the present invention considered the facts that the shear rate of a resin in injection molding differs greatly from that in extrusion molding, the shearing rate being in the range of from 10000 to 100000 sec.sup.-1 in injection molding and being in the range of from 100 to 1000 sec.sup.-1 in extrusion molding, and that the nonflammable resin to which the moldability-improving agent is added is not a resin of ordinary molecular weight but is a special nonflammable injection-molding resin which comprises a styrene type resin having a high content of low-molecular weight .alpha.-methylstyrene and containing a post-chlorinated vinyl chloride type resin of low polymerization degree. Based on the above consideration, the present inventors comprehensively investigated moldability-improving agents of the alkyl acrylate resin type to find modifier which does not lower the effect of the moldability-improving agent, decrease impact strength, nor cause peeling of molded articles. As a result, it has been found that an alkyl acrylate copolymer having a specific composition achieves the above objective."}
{"text": "Enterprises adopt workflow management tools to automate their business processes and handle complex batch processing. Batch applications are computationally data intensive and if executed during business hours by an enterprise, may impact its critical customer front-end applications due to their high resource cost. In addition, delay or failures in batch-application execution have significant business impact potential due to their high do-over cost in time and compute resources.\nTo make quick business decisions, real-time inputs to business intelligence applications are critical and these applications have to perform complex operations and generate reports with speed and accuracy. Delays may impact business decisions and produce a rippling effect on business operations. In addition, front-end applications such as e-commerce and financial analysis platforms need to be served a large set of data in real-time as they perform complex and computationally intensive algorithms based on customer requests. Therefore, in a time when enterprises are experiencing pressure to make their business more agile, their business applications need to perform and scale optimally and, preferably, in real-time.\nCurrent batch application processing solutions are based on dedicated clustering technologies for high availability and performance. Although clustering technologies allow computing infrastructure to scale-out on-demand to a certain extent they may be unable to grow beyond their local spatial environment and are generally homogenous computing clusters. However, batch processing applications may require a highly heterogeneous and distributed environment for fast real-time scale out that integrates fault tolerance. Therefore, merely clustered based batch solutions may be unable to meet an enterprise's demands."}
{"text": "1. Field of the Invention\nThe invention relates to a method of encoding/decoding effectively a fading image and a dissolve image and an apparatus therefor.\n2. Description of the Related Art\nIn a video encoding standard scheme such as ITU-TH.261, H.263, ISO/IECMPEG-2, and MPEG-4, a motion compensated predictive interframe encoding is used as one of encoding modes. As a predictive model in the motion compensated predictive interframe encoding is adopted a model indicating the most high predictive efficiency when luminosity does not vary in a time axis. In the case of the fading image that luminosity of the image varies, for example, in the case of fading-in from a black image to a normal image, a method of performing a prediction adequately according to a change of luminosity of the image is not known. Consequently, there is a problem to need the large number of encoded bits for the purpose of maintaining an image quality in a fading image.\nFor example, U.S. Pat. No. 3,166,716 discloses a technique of coping with the problem by detecting a fading image area and changing an allocation of the number of encoded bits thereto. Concretely, in the case of the fade-out image, the large number of encoded bits are allocated to an beginning part of the fade-out that brightness varies. The last part of the fade-out usually reduces allocation of the number of encoded bits since it normally becomes a monochrome image resulting in making the encoding easy. With such a way, a total image quality is improved without increasing the total number of encoded bits.\nOn the other hand, U.S. Pat. No. 2,938,412 discloses an encoding system that deals with the above problem in a fading image by compensating for a reference image according to two parameters of an amount of luminance change and an amount of contrast change.\nThomas Wiegand and Berand Girod, “Multi-frame motion-compensated prediction for video transmission”, Kluwer Academic Publishers 2001, provides an encoding system based on a plurality of frame buffers. This system intends to improve a predictive efficiency by generating a predictive image selectively from a plurality of reference frames saved in the frame buffers.\nHowever, the system of U.S. Pat. No. 3,166,716 improves an image quality without increasing the total number of encoded bits in encoding a fading image by detecting a fading image area, and changing an allocation of the number of encoded bits. For this reason, there is the advantage that the encoding can be realized within the framewark of an existing encoding system. However, since the prediction efficiency is not essentially improved, a notable improvement of the encoding efficiency cannot be expected.\nOn the other hand, the system of U.S. Pat. No. 2,938,412 has a merit that the predictive efficiency on a fading image improves. However, the predictive efficiency to be enough for so-called a dissolve image (referred to as a cross fade image) which an image gradually varies from an image to another image cannot be obtained.\nThe system of Thomas Wiegand and Berand Girod cannot sufficiently deal with a fading image and a dissolve image, and cannot improve a predictive efficiency even if a plurality of reference frames are prepared.\nAccording to the prior art as described above, the large number of encoded bits are required for the fading image and dissolve image to be encoded with high image quality. There is a problem that improvement of the encoding efficiency cannot be expected."}
{"text": "The present embodiments relate to video analytics. In particular, configuration of multiple cameras is provided for video analytics.\nIn order to meet the increasing needs of the security and safety, surveillance systems are becoming more complex. Besides the basic capabilities of video surveillance, such as recording and playback, video analytics technologies are used to enrich the functionality. A person or vehicle may be detected and tracked, events may be detected, an alarm may be triggered, and/or forensic searching may be provided.\nIn order to achieve accurate performance in terms of detection and false alarm rate, every camera needs to be configured properly. The camera configuration task varies based on the different functionalities. Camera configuration is time consuming. For example, a tunnel monitoring system monitors traffic flow and generates events, such as slow vehicle, stopped vehicle, traffic jam, wrong directional driver, and pedestrian detection. An experienced worker may need about eight minutes to configure a camera, during which the worker may click the mouse buttons more than 100 times. With forty cameras for every mile inside the tunnel, 5.5 hours and 4000 mouse clicks are needed to configure the cameras for every mile.\nFor indoor people detection and tracking functionalities, configuration tools are typically provided for a single camera. As a person travels from the field of view of one camera to the field of view of another camera, the person is re-identified. The facial recognition or other detection is performed again. To configure a typical site with 40 cameras, a worker may click 120 times to define camera positions and more than 150-600 times overall to define the calibration planes for all of the cameras."}
{"text": "A well contact diffusion layer (a so-called “guard ring”) is frequently used in a semiconductor device to electrically connect semiconductor elements such as a transistor to a well diffusion layer. Contacts connecting wirings from a semiconductor device to a guard ring are often formed to be circular or elongated-circular in a planar shape.\nIn terms of reducing a chip size of a semiconductor device, contacts of a guard ring are desirably circular-hole contacts having a circular shape. However, if all the contacts of a guard ring are circular-hole contacts, contact resistances to wells are increased when some of the circular-hole contacts are not opened.\nOn the other hand, in terms of reducing contact resistances to wells, contacts of a guard ring are desirably elongated-hole contacts having an elongated shape. However, if all the contacts of a guard ring are elongated-hole contacts, the chip size may be increased. For example, the width of the guard ring needs to be increased to match the major axis of the elongated-hole contacts. Alternatively, distances between adjacent wirings need to be increased to match the major axis of the elongated-hole contacts."}
{"text": "This invention is directed generally to harvester reels and specifically to a cam operated rotating reel with improved means thereon for selective lateral adjustment of the tine bars relative to their support bearings.\nCertain types of crop harvesters employ a rotating reel at the forward end of the header to feed material into a mower and to rake the severed material rearwardly into conditioning rolls, or other crop manipulating mechanisms. The reel generally includes crop engaging fingers, or tines, affixed to elongate circumferentially spaced tine bars which are journalled in two or more laterally spaced support members fixed to the reel axis. Usually, though not necessarily always, the support members are spiders, with the tine bars journalled in each leg thereof. The attitude of the tines is controlled by a cam means which rocks the tine bars through a predetermined angle as the reel rotates.\nThe type of reel to which the instant invention is most applicable is well known in the prior art as exemplified in U.S. Pat. No. 3,472,008. Though intermediate supports may be used as required, the drawings of this patent show two laterally spaced spiders fixed to a common axis. Each leg of each spider terminates radially with a horizontally disposed bearing held in position by a bearing housing. Horizontal elongate tine bars extend through the aligned bearings in the spiders and are freely rotatable therein. Each tine bar carries a multiplicity of horizontally spaced tines fixed thereto for engagement with the crop material. An arm and cam follower is fixed to one end of each tine bar to be engaged with a channel-shaped closed loop cam on the interior surface of the reel housing.\nAs is common in reels of the type under consideration, the cam is constructed such that the tine bars, and thus the tines, rotate less than 360.degree. relative to the support bearings as the reel rotates 360.degree.. This relative rotation is advantageous in that it pulls the crops across the cutter bar in a continuous manner and provides improved feeding characteristics into the auger, or the nip of the conditioning rollers.\nThe cam structure, as readily seen in above-mentioned patent, comprises horizontal opposing camming surfaces forming a side opening into which the cam follower, or roller, is inserted. It should be readily apparent that integration of the reel structure into the housing is facilitated by some means for lateral adjustment of the tine bars relative to the support bearings. More specifically, during assembly lateral adjustment of the tine bars is necessary to ensure proper alignment between the cam follower and the cam surfaces. The means employed to make this adjustment must not interfere with the rotational characteristics of the tine bar relative to the bearings, but yet must be capable of selectively rigidly fixing its lateral location.\nThe prior art has employed various devices to accomplish the above-described adjustability. For example, U.S. Pat. No. 3,641,750 shows a pair of washer-like collars located on the tine bar and fixed thereto by set screws. Furthermore, each collar has one surface thereof engaged with a support bearing housing, thus maintaining a fixed distance between spiders. Besides the multiplicity of elements and the increased likelihood of misadjustment inherent therewith, such an arrangement adds an undesirable degree of complexity to the assembly operation.\nAnother prior art method of making the adjustment is to use a clamp which straddles the tine bar support bearing and the end of the spider and grips the bar on both sides thereof. The tine bar is laterally adjusted by loosening the clamp and sliding the bar endwise to its desired position and then retightening the clamp. Though effective, the adjustment clamp is exposed to crop materials and debris, and is thus potentially easily jammed. If the clamp cannot rotate about the end of the spider, the tine bar will also not be able to rotate and function properly.\nOther methods employed in the prior art use various combinations of washers and cotter pins to, at best, provide inconvenient and awkward adjustment of the tine bar. The invention to be described in detail below is to a unique adjusting device which exhibits none of the shortcomings of the known prior art."}
{"text": "Conveyors are used in a number of applications to convey articles from a first point to a second point. Some common conveyor types include roller conveyors, which are typically comprised of a plurality of motorized and non-motorized rollers, and belt conveyors, which include belts that are driven by one or more motorized rollers. The conveyor may be controlled by a single centralized controller, or may be divided into multiple sections, or zones that are each controlled by a separate controller. In either case, the one or more controllers generally control the conveyor system by selectively activating motorized rollers or motors that operate the belt of the conveyor to move an article along the conveyor. In conventional conveyor controllers, the controller may also determine whether an article is being conveyed based on signals from presence sensors configured to detect the presence of an article on the conveyor, such as photo-eye sensors.\nMotorized Drive Roller (MDR) conveyor systems use conveyor rollers that are driven by electric motors to move articles along the conveyor route. These types of conveyor systems typically include multiple independently controlled zones, with each zone having a single motor and an article presence sensor operatively coupled to a zone controller. The rollers of a zone are typically coupled to the driven roller so that the rollers operate in concert to selectively move an article through the zone in response to activation signals from the zone controller. The zone controller may be operatively coupled to adjacent zone controllers so that activation of the zone can be coordinated with the movement of articles upstream and downstream of the zone. A conveyor system may be implemented by arranging a plurality of these zones sequentially so that articles are transported in the conveying direction. These zones may include straight sections, curved sections, diverting sections, and merging sections so that articles may be received from and discharge to multiple locations. The modular characteristics of MDR conveyor zones facilitate reconfiguring conveyor systems by adding, removing, and reconfiguring zones within an existing conveyor system. An existing MDR conveyor system can thereby be changed and expanded as the needs of the conveyor operator evolve.\nTo achieve efficient conveyance of articles, a conveyor system should convey articles with short intervals, or gaps, between articles. Short spacing intervals increase article traffic density, and thus conveyor system throughput for a given conveyor operating speed. Short spacing intervals also reduce the overall length required in accumulation conveyor systems for a given article accumulation capacity. However, in order to prevent articles from colliding and/or being delivered to the wrong destination, conventional modular conveyor systems limit the number of articles to one per zone. Achieving short spacing intervals in a modular conveyor system using conventional zone controllers therefore requires short zones. This short zone requirement can significantly drive up costs of conveyor systems that have long spans or that accumulate large numbers of articles due to the large number of zones required to span the length of the conveyor system.\nConsequently, there is a continuing need for improved methods and systems for modular conveyor control that reduce the cost of implementing modular conveyor systems, and that maintain or increase efficiency and throughput of modular conveyor systems."}
{"text": "1. Technical Field\nThe present invention relates broadly to improvements in tactile sensing of acoustic and audio signal for hearing impaired persons. The invention has particular, although not exclusive, utility in enhancing the ability of hearing impaired persons to enjoy the audio portion of television signals.\n2. Discussion of the Prior Art\nTactile devices for the deaf have been studied for approximately sixty years, and within the last eight years have been produced as commercial items. The essential concept embodied in all of these devices is to convert sound signals into vibrations that can be felt on the skin via the tactile sense. Naturally enough, the greatest emphasis has been on the ability to recognize speech sounds, although the practical value of recognizing various background sounds, such as bells, footsteps and so on, are also well appreciated. Since the tactile sensory system is unable to perceive vibrations much above 800 to 1000 Hz, and since significant input from signals to at least 4000 Hz is required to impart important speech data, a significant portion of the research and development work in this field has addressed methods for processing sound signals such that information from the higher frequencies can be transmitted to a user. Virtually all work in the field has focused on wearable devices which can be used on a more or less continuous basis; accordingly, a further constraint has been that the methods chosen result in small, cosmetically acceptable designs. Using these guidelines, the kinds of systems that have evolved all use small skin transducers, worn either on the wrist, chest or around the back of the neck. While it is well known that placing the transducers somewhere on a hand yields very good perceptual results, these locations are avoided for reasons relating to utility. Typically these transducers are resonant at a single frequency, usually 250 Hz as this is the most sensitive frequency for the tactile sense, and processing methods are employed that encode the desired signals at this rate. In the more advanced systems a number of such transducers are used, each being assigned to a different portion of the acoustic frequency band and worn at a respective position on the skin. Thus, in these systems, the incoming sound signal is separated by a processor into spectral segments, each segment being represented at a respective location on the skin as a 250 Hz signal with a varying time envelop and intensity according to the spectral content of the original signal.\nThe most advanced of these so-called \"multiple channel\" tactile devices currently available use seven such transducers arranged side by side with channels divided to cover the sound frequency range from 200 Hz to 7000 Hz, divisions being selected according to data on vowel frequencies commonly known in the speech discipline. The processing arrangement permits no more than two transducers to be actuated at the same time. It has been found that this method is desirable both to provide clarity of perception and to conserve battery life.\nThe reason for using resonant transducers is that they are more efficient than wideband transducers, at least as regards presently available technology. This, in turn, allows compact battery supplies to provide reasonably long operation with attendant cosmetic acceptability, but only at the cost described below.\nIt has been recognized that wideband transducers allow more easily perceived patterns with richer information content to be applied to the skin. In particular there are certain sound information components, notably voicing data in speech and attack waveform data in music, that are better represented by direct presentation of the unprocessed signals for the lower frequency range of the incident sound signal. This is true for that portion of these signals lying between about 50 Hz to 800 Hz; frequencies above this tend to be so dimly perceived by the tactile sense that there is little if any information communicated without some kind of encoding being used, such as the transducer location encoding described above.\nIt may be that improved perception could be obtained with a hybrid system using a similar multiple channel scheme wherein the lower frequency channels use no encoding, being excited directly by the lower spectral portion of the original sound signal, and only the upper band spectral components are encoded as described above. However, the technical methods for obtaining such a system in wearable form are not currently available in a configuration that also meets cosmetic requirements.\nOn the other hand, if one relaxes the cosmetic requirements, as would be appropriate for the intended uses of the present invention, then broadband transducers such as loudspeakers or small motors, as described hereinbelow, provide a means for supplying the desired direct tactile representations of lower frequency sounds.\nIn order to provide a frame of reference, it is desirable to note the specifics of the narrowband transducers currently used in the tactile field. A typical unit is approximately one inch long by 0.5 inch thick by 0.7 inch wide and weighs about seven grams. Its low weight and flat configuration is appropriate for the usual methods of mounting. Since all signals applied to such a transducer are encoded at a single frequency, typically 250 Hz, requirements for bandwidth are only related to the required response time. A well designed unit follows waveform changes at better than a 20 msec rate, adequate for virtually all speech sounds and all environmental sounds of interest. For such a device a typical peak driving energy to obtain maximum perceptual feel is on the order of 350 mw. It has been found over a period of years of actual use that the average duty cycle for these kinds of systems is on the order of 0.2. Accordingly, for a wearable system using one such transducer, the expended energy/hour is on the average about 70 mwatt-hours. A typical rechargeable battery pack used to power these devices, constrained in size according to the usual cosmetic considerations, stores about 1600 mwatt-hours. Hence, as is verified in many thousands of hours of actual use, a system using one such transducer will operate, on the average, for about twenty-three hours between charges. Generally, the decrease in hours of operation between charges does not follow linearly with an increasing number of transducers because, depending on design parameters, it is not true that all transducers in a system have the same duty cycle. For the seven channel device mentioned previously, where only two transducers at most are activated simultaneously, the time of operation between charges is approximately 15 hours.\nIn contrast, a typical mass-produced small motor requires about 1.5 watts in the application mode described herein to elicit the same perceptual intensity; it weighs about forty grams and has a cylindrical configuration typically on the order of one inch in diameter and one inch in length. It is clear that neither the configuration nor the power requirements lend themselves to application in a wearable form and that the weight also mitigates against this use. However, for a hand-held or hand-contact application in which larger size is allowable and hence a larger space can be allocated for the battery pack, the use of such small motors is not precluded. If loudspeakers are utilized, they require only 350 mw to achieve the same level drive for the applications described herein; however the loudspeakers are typically even larger and heavier than the small motors. Again, for hand-held or hand-contact applications this larger size and weight is not objectionable.\nFrom the foregoing it will be appreciated that the present invention is intended to teach the use of small mass-produced motors or loudspeakers in tactile transducers for the deaf and hearing impaired in those applications where large size and low efficiencies, as compared to the usual tactile transducers, are not serious detriments. In particular, but not exclusively, this would apply to transducing television audio signals whether or not accompanied by captioning.\nThe present invention has particular, but not exclusive, utility in conjunction with television programming incorporating closed captioning or similar techniques employed to provide titling on the video screen to represent spoken words in the audio signal. The closed captioning titling technique is described and defined in Report No. E-7709-C, revised May 1980, by the Public Broadcasting Service. The disclosure in that report is expressly incorporated herein. Typically, a closed captioning decoder processes the closed captioning signal that is included in line twenty-one of the NTSC video signal. Decoding results in titling that is overlayed on the television screen, the titling content corresponding to speech in the audio portion of the received television signal. The present invention, by providing a low cost wideband transducer that can be used in conjunction with closed captioning, not only enhances appreciation of television programs for hearing impaired persons; it also permits such persons to be trained in the use of the transducers since the speech signals being tactually transduced are simultaneously viewed as titling on the television screen."}
{"text": "The present disclosure relates generally to digital imaging devices and, more particularly, to systems and method for processing image data obtained using an image sensor of a digital imaging device.\nThis section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present techniques, which are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present disclosure. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nIn recent years, digital imaging devices have become increasing popular due, at least in part, to such devices becoming more and more affordable for the average consumer. Further, in addition to a number of stand-alone digital cameras currently available on the market, it is not uncommon for digital imaging devices to be integrated as part of another electronic device, such as a desktop or notebook computer, a cellular phone, or a portable media player.\nTo acquire image data, most digital imaging devices include an image sensor that provides a number of light-detecting elements (e.g., photodetectors) configured to convert light detected by the image sensor into an electrical signal. An image sensor may also include a color filter array that filters light captured by the image sensor to capture color information. The image data captured by the image sensor may then be processed by an image processing pipeline, which may apply a number of various image processing operations to the image data to generate a full color image that may be displayed for viewing on a display device, such as a monitor.\nWhile conventional image processing techniques generally aim to produce a viewable image that is both objectively and subjectively pleasing to a viewer, such conventional techniques may not adequately address errors and/or distortions in the image data introduced by the imaging device and/or the image sensor. For instance, defective pixels on the image sensor, which may be due to manufacturing defects or operational failure, may fail to sense light levels accurately and, if not corrected, may manifest as artifacts appearing in the resulting processed image. Additionally, light intensity fall-off at the edges of the image sensor, which may be due to imperfections in the manufacture of the lens, may adversely affect characterization measurements and may result in an image in which the overall light intensity is non-uniform. The image processing pipeline may also perform one or more processes to sharpen the image. Conventional sharpening techniques, however, may not adequately account for existing noise in the image signal, or may be unable to distinguish the noise from edges and textured areas in the image. In such instances, conventional sharpening techniques may actually increase the appearance of noise in the image, which is generally undesirable. Further, various additional image processing steps, some of which may rely on image statistics collected by a statistics collection engine, may also be performed.\nAnother image processing operation that may be applied to the image data captured by the image sensor is a demosaicing operation. Because the color filter array generally provides color data at one wavelength per sensor pixel, a full set of color data is generally interpolated for each color channel in order to reproduce a full color image (e.g., RGB image). Conventional demosaicing techniques generally interpolate values for the missing color data in a horizontal or a vertical direction, generally depending on some type of fixed threshold. However, such conventional demosaicing techniques may not adequately account for the locations and direction of edges within the image, which may result in edge artifacts, such as aliasing, checkerboard artifacts, or rainbow artifacts, being introduced into the full color image, particularly along diagonal edges within the image.\nAccordingly, various considerations should be addressed when processing a digital image obtained with a digital camera or other imaging device in order to improve the appearance of the resulting image. In particular, certain aspects of the disclosure below may address one or more of the drawbacks briefly mentioned above."}
{"text": "1. Field of the Invention\nThe present invention is directed to a TONG BEARING and in particular, to a machinable TONG BEARING that may be reconditioned for multiple uses and an extended useful life.\n2. Description of the Prior Art\nTong systems are utilized in the oil well industry for attaching to and tightening various types of rods and tubing. Such tongs are utilized with a backup system that holds one element while a second element is gripped by the tong and rotated to connect and disconnect the two elements.\nA typical tong includes a housing around a gripping portion of the tong, with a gripping portion being driven and rotating relative to the housing. Conventional tongs utilize a series of bearing elements to facilitate rotational movement with less friction and wear. Such rollers, ball bearings or other conventional bearings are subject to the harsh conditions encountered at a well or drilling site. Therefore, such single use conventional bearings are subject to wear and/or failure, requiring frequent replacement. Moreover, in order to ensure smooth and continuous operation, multiple bearings are required at multiple positions spaced apart on both the top and bottom of the tong housing around the gripping assembly. The maintenance and replacement of the large number of bearings is both time consuming and expensive. Moreover, the maintenance required and high rate of failure can cause extended down time for the tong, affecting reliability, causing additional delays and adding to operational costs.\nIt can be seen then that a new and improved bearing system is needed for tongs. Such a bearing system should eliminate the large number of bearings and small parts currently needed. Moreover, the bearing system should provide for greater reliability and less down time. Such a bearing should be easily removed and interchanged with other bearings. Moreover, the utility of such a bearing is improved if a bearing element can be reconditioned and used again. The present invention addresses these as well as other problems associated with tong bearings."}
{"text": "In computer graphics, traditionally colors are represented using a combination of primary colors. For example, in the RGB color space, the colors Red, Green and Blue are blended together to get a range of colors. In order to manipulate images in the RGB color space, scripting languages are designed to allow the users to create scripts that describe image-processing operations at a high level. For example, a script can be written to handle transition between two video clips or between one video clip and one still video. The scripts do not necessarily work only with transitions. They may also be written to work with single source clips. For example, a script that takes an existing clip of video and drop colors on it. The scripts may also be “generators” which have no inputs but create an output.\nFinal Cut Pro (FCP) is a movie editing and creating software produced by Apple Computer, Inc. of Cupertino, Calif. (“Apple”). It includes a scripting engine. The FCP scripting engine allows users to write scripts which perform various image manipulations, ranging from basic operations such as “blend” “channelfill” and “multiplyChannels” to more complex operations such as “levelmap” and “colorkey.” Of course, most scripts combine more than one of these operations to build interesting effects or transitions.\nThese functions operate on one or more images. Consider these examples. “Channelfill” fills in one or more color channels of a destination buffer with the color values passed to it. “Levelmap” applies a lookup-table to one or more channels of “source” and puts the results of the lookup into “destination”. “Blend” blends “source 1” and “source 2” together into “destination”, based on a ratio passed in. These scripting engine commands, or sometimes referred to as image processing calls or functions, were shipped with many other commands in version 1.0 of FCP from Apple. In addition, there are many scripts written by the users of FCP.\nIn the past, it was assumed that the scripts were written to operate in the RGB color space, 8 bits per color component, and the scripting engine only performed the image processing in the RGB color space."}
{"text": "There are different solutions for the production of energy from sea waves, for example, solutions based on wave height, wave movement, etc. The present invention utilizes both the wave length and the wave width by the forward wave motion, the backward wave motion and the force of gravity in the most efficient way, something that other proposed solutions do not.\nThe following are examples of the innovation inherent in the present invention as compared with the patents enumerated below, both from the point of view of the operating principle and structure, and from the point of view of energy considerations—allowing greater energy to be attained and produced.    1. (i) U.S. Pat. No. 5,084,630 utilizes only the width of the wave and not the length thereof.            In this patent there is disclosed a system consisting of number of units provided with paddles for each section of wave width. At the bottom of each paddle there are plates. The plates are spaced laterally one from another relative to the direction of wave movement (FIG. 1, the patent top view, shows it). Each paddle activates a pump. In this way, the wave power is not sufficiently utilized because after the wave strikes the paddle and leaves it, it continues to advance and is not longer in contact with the paddle, In FIG. 1, part of the 30 hits the plate 44, leaves the paddle and the wave's part continues to move without utilized its power any more. Other parts of the wave 30 hit the plates 46, 48, 50, 52, 54 and 56, each one wave's part hits only one plate and continue to move without utilizing its power any more. Hence, the wave's power is not fully utilized. The system is disclosed in U.S. Pat. No. 5,084,630 utilizes variations in sea level. For this purpose it is provided with a dedicated subsystem, including a set of separates pumps. Due to this provision the whole system is complicated, as it requires additional control means for coordination of the paddle disposition with the variation in sea level. Furthermore, as the disposition of the paddle is continuously changed and the paddle is not submerged deeply enough to maximally utilize available energy, it is not possible to utilize the energy of those wave's components that move under the wave.        \nThe present invention utilizes:                most efficiently both the length and the width of the wave in the following way:        The lower extremities of the paddles are parallel to the wave width. They are arranged one after other in such a way that from a top view, they are arranged in series. Then, when the wave's part strikes the lower extremity of the first paddle it continues to move, and also strikes the other paddles, until the wave's part is flattened. In this way, the power of the wave is consumed more efficiently so as to get maximum energy from the wave. The energy of the next wave coming after the first wave will be utilized in a similar way.        \n(ii) U.S. Pat. No. 5,084,630 utilizes:                Only the energy in one direction of each paddle. When a paddle, 66+68+70+82+84 as it is shown in FIGS. 3-7, moves pivotal forward, the piston 124+128 as it is shown in FIG. 8, pushes oil within hydraulic cylinder 130 to give the energy to the generator 196, through 188, 146 and accumulartor 140. The oil comes from hydraulic tank 208 through pipe 210. When the paddle returns backward after the wave leaves the paddle, the piston 124+128 returns back and presses some air and/or oil, if desired back to the tank 208 through inlet pipe 212. The accumulator 140 does not get oil and the generator 196 does not get energy. The U.S. Pat. No. 5,084,630 has a one direction piston. The patent uses the paddle energy in one direction, the patent does not utilize the paddle's energy.        \nThe present invention utilizes:                The whole paddle energy. The invention has for each paddle, a bi-directional piston. When a paddle move forward with the wave, the piston inside the hydraulic cylinder, pushes oil to the first accumulator, after the wave leave the paddle, it returns backward and the same piston pushes oil to the second accumulator. The two accumulators send the oil throgh pipes to an hydraulic engine that rotates a generator. The present invention receive energy from both paddle pivotal forward motion and paddle pivotal backward motion.            2. (i) PCT/IL99/00258 or AU PP8932 or GB9800394.0 patent requests (they are the same requests), (hereinafter: “PCT/IL99/00258 patent requests”) show in the application—page 5 paragraphs 4,6, and 7, that the invention has one rod with two pistons with two hydraulic cylinders, both of them send the oil to one accumulator. It is theoretical, in reality the only way that an accumulator can be filled by hydraulic cylinders is, one by one. The hydraulic cylinders above the paddle rod operate one by one and the hydraulic cylinders below the paddle rod operate one by one also, but all the hydraulic cylinders, above and below paddle rod, don't operate one by one. In this way the accumulator can't be filled and work properly. The PCT/IL99/00258 patent requests can be used only to produce energy up to the accumulator. The accumulator should send oil to hydraulic engine that operates the generator. The PCT/IL99/00258 patent requests can't make progress from the accumulator to the generator, because the accumulator does not work. The PCT/IL99/00258 patent requests not can work and produce energy because it uses one accumulator.\nIn the present invention, there are two accumulators (it is shown below in FIG. 4). All pistons oil exit in the first direction, fill the first accumulator with oil one by one, when the paddles move forward, and all pistons oil exit in the second direction, fill the second accumulator with oil one by one when the paddles move backward. This invention can be used to produce electric energy with two accumulators that send the oil to hydraulic engine that rotates the generator.\n(ii) PCT/IL99/00258 patent request claims don't say that in backward paddle unit motion the energy is utilized, i.e. after a wave leave the paddle's plate and the paddle falls pivoting with its gravitation and the backward wave's power, this paddle energy is not utilized. The only energy that is utilized is the energy of a wave that hits the paddle's plates. The gravitation energy of each paddle and the backward wave's energy aren't utilized.\nThe present invention shows in claim 1, that either, in forward paddle unit motion, the paddle's energy is accepted and consumed and, in backward paddle unit motion, the paddle's energy is accepted and consumed.\nThe present invention utilizes:                The whole paddle energy. The invention has a bi-directional piston for each paddle. When a paddle move forward with the wave, the piston inside the hydraulic cylinder, pushes oil to the first accumulator, after the wave leaves the paddle, it returns backward and the same piston pushes oil to the second accumulator. The two accumulators send the oil through pipes to an hydraulic engine that rotates a generator. The present invention receives energy from both motions: paddle pivotal forward motion and paddle pivotal backward motion.        \nThe present invention produces energy in the most efficient way, and the invention's technology is better than PCT/IL99/00258.    3 (i) PCT/IL01/00271 patent request shows in their application different method and different drawings than the present invention, as explained in the PCT/IL01/00271 International Preliminary Examination Report patent request, Item I, in which the FIG. 1-FIG. 6 are the present invention.\nThe following is Item I in the above PCT/IL01/00271 report:                “The present examination authority considers that the amendments files with the letter dated Jun. 17, 2002 (received on Oct. 21, 2002) introduce subject-matter which extends beyond the content of the application as originally filed, contrary to Article 34(2)(b) PCT. The amendments concerned are the following:        \nWith the above-mentioned letter the application has amended the drawing FIGS. 1 and 2 on drawings page 1/6, and has introduced the drawings FIG. 5 on drawings page 4/6 and drawing FIG. 6 on drawing page 5/6. The amended “connections” in FIGS. 1 and 2, as well the two additional drawings FIGS. 5 and 6 where not contained in the drawings pages 1/5-5/5 as originally filed.\nIt is not apparent how the amended “connections” in FIGS. 1 and 2, or the additional FIGS. 5 and 6 are directly and unambiguously derivable from the content of the application as originally filed. Although the respective structural arrangement shown in said FIGS. 1, 2, 5, and 6 may indeed not be excluded from the original subject-matter, it remains that it is not explicitly indicated in the application as originally filed that said respective structural arrangement of FIGS. 5 and 6 necessarily has to be realized in order to perform the claimed invention, or how/why the original “side-to-side connection” of FIGS. 1 and 2 may be superseded by a “bottom-to-bottom connection”.\nHence, the amendments in FIGS. 1 and 2, as well as the additional FIGS. 5 and 6 implicitly introduce subject-matter which extends beyond the content of the application as originally filed, contrary to the requirements of Article 34(2)(b) PCT.\nAs foreseen by rule 70.2(c) PCT, this international preliminary examination report is being established as if the above mentioned amendments and additional had not been made. Hence, it is theoretically assumed for the rest of this report that the drawings figures on files are those shown on drawing page 1/5-5/5 as originally filed. Similarly, the description parts related to the amendments and additions in the drawings are considers not to have made.\nForm PCT/Separate Sheet/409 (Sheet 1) (EPO-April 1997)”\nWhile the connection of the rod 3 is to the side of 14 as in PCT/IL01/00271 patent request and a wave arrives and hits there is an energy loss that comes from the momentum of the 14 side distance multiply by the paddle's weight. When the plant is large each of the paddles have a big weight, which bring to a high momentum with high energy loss.\nFollowing the PCT/IL01/00271 International Preliminary Examination Report, I was forces to ask the PCT patent request for the present invention with the new FIGS. 1 and 2 that I sent to the PCT/IL01/00271. The figures were not published because the PCT/IL01/00271 International Preliminary Examination Report was not published. In addition: please note to the 3rd paragraph above this paragraph, from the words: “As foreseen by rule 70.2(c) PCT” till the words: “the drawings are considers not to have made.”\nThis invention bottom-to-bottom connection brings an essential change with high advantage, that the plant energy production is higher than PCT/IL01/00271 patent request.\n(ii) The PCT/IL01/00271 patent request uses, in each one of the paddle unit, as it is written in its claim 1, only one crankshaft mechanism and it can be applied only to very small waves. It can be applied for high waves. When a high wave, (e.g. 5 meters wave height, etc. ) arrives, it hits the paddle plates with high force and one crankshaft mechanism, which is the piston 1 in the present invention drawings, should be very large and the piston size is not realistic piston to be operated.\nThe piston's area is calculated according to the following formula:                The piston area=(The wave's force)/(The pressure in the hydraulic system). Since the piston size is so large and not realistic, the plant must have at least two piston to that each one of the two (for example two) pistons' area is half size of the original one piston and only by splitting the one piston to at least two pistons, the plant can work. The present invention using at least two crankshaft mechanisms, as it is written in the present invention claim 1, which can be applied for high waves that bring high power for high plant electric energy production.        \nEach one of the above two alterations in respect of PCT/IL01/00271 patent request is an essential alteration. According the above 1st alteration the, the present invention plant energy production is much more than in PCT/IL01/00271 patent request, and according to the 2nd alteration, the present invention can work and produce energy with high and low waves, while PCT/IL01/00271 patent request can work only with low waves.\nThe present invention produces energy in the most efficient way, and the invention's technology is better than PCT/IL011/00271 patent request.    4. (i) Patent GB 384603A and Patent JP 57081168 have one crankshaft mechanism for the whole paddles.\nThe energy consideration is:                The power of each wave is known on the basis of the wave's height and width. The wave strikes the first paddle, the second paddle, the third paddle and so forth, until the wave completely flattens, and then the next wave comes along. Just like an accumulator that is depleted, so the system takes all the power minus losses.        \nIn General:\n  Power  =            Energy      Time        =                  The        ⁢                                  ⁢                  plant          '                ⁢        s        ⁢                                  ⁢        power            =                                                                  The                ⁢                                                                  ⁢                total                ⁢                                                                  ⁢                energies                ⁢                                                                  ⁢                of                ⁢                                                                  ⁢                all                ⁢                                                                  ⁢                paddles                ⁢                                                                  ⁢                situated                                                                                        in                ⁢                                                                  ⁢                series                ⁢                                                                  ⁢                between                ⁢                                                                  ⁢                two                ⁢                                                                  ⁢                waves                                                              Time          ⁢                                          ⁢          between          ⁢                                          ⁢          two          ⁢                                          ⁢          waves                    \nThe mathematical principle of calculating the wattage:\nP—The power as calculated according to its height and breadth.\nμ—Paddle efficiency.\nNumber ofWave power before the NthPaddles - NPaddlePaddle Power1PP * μ2P − P * μ(P − P * μ) * μ3(P − P * μ) − (P − P * μ) * μ =P(1 − μ)2 * μP(1 − μ)2.........P - The power as calculated according to its height and breadth.μ - Paddle efficiency.\nThus, the series continues, so that the greater the number of paddles, the greater the system power, up to the maximum P.\nThe rod that links the last paddle to the crankshaft mechanism in Patent GB 384603A and Patent JP 57081168 is extremely big and long because of the long distance between two waves.\nIn order to obtain maximum power, it is necessary to provide at least 15 paddles. When the waves are 2 meters to 5 meters height, the distance between two waves is 50-100 meters, so that said rod must also be 50-100 meters!\nBecause such a great rod length is required, the profile of the rod proposed in GB 384603A or in Patent JP 57081168 is extremely thick and large, so that it should not bend under the force of the waves, and quite impractical.\nBecause of the weight of the rod proposed in GB 384603A or in Patent JP 57081168, and the slow wave speed prevailing under said conditions, it is not possible to attain the minimum speed of 500 rpm required to rotate the flywheel of the oil pump (linked to the accumulator that drives an oil engine, which in turn drives a generator and produces the electric energy).\nThus, it should be emphasized, that the arrangements proposed in GB 384603A or in Patent JP 57081168 are completely impractical where waves higher than 2 meters are concerned.\nThus, Patent GB 384603A and Patent JP 57081168 are only applicable when Waves are small and the distance between two waves is short. Where higher waves are concerned, the entire plant proposed by GB 384603 A or by Patent JP 57081168 become dysfunctional and cannot be used.\nThe present invention has a crankshaft mechanism for each paddle unit, without used a rod to link among the paddles. Each paddle unit is operated independently.\nThe present invention is equally practicable for both small and large waves and can be used to obtain high power capacities under both sea conditions, and this is the technological innovation of my patent application over the patents: GB 384603A and JP 57081168.\nThe present invention shows below in claim 1 (the main claim ), that it incorporates the features referring to crankshaft mechanism and a piston and also to providing each of the paddle units with a dedicated crankshaft mechanism to show the difference in patents GB 384603A and JP 57081168 in which the whole plate has one crankshaft mechanism.\n(ii) Patent GB 384603A and Patent JP 57081168 don't utilize the energy of backward paddle unit motion, after a wave leaves the paddle's plate and the paddle falls pivoting with its gravitation and the backward wave's power, this paddle energy is not utilized. The only energy that is utilized is the energy of a wave that hits the paddle's plates. The gravitation energy of each paddle and the backward wave's power, aren't utilized.\nThe present invention utilizes the whole paddle energy.                The invention has a Two-Direction piston for each paddle. When a paddle moves forward with the wave, the piston inside the hydraulic cylinder pushes oil to the first accumulator, after the wave leaves the paddle, it returns backward and the same piston pushes oil to the second accumulator. The two accumulators send the oil through pipes to an hydraulic engine that rotates a generator. The present invention receives energy from both motions: paddle pivotal forward motion and paddle pivotal backward motion.        \nThe present invention shows in claim 1 that either, in forward paddle unit motion, the paddle's energy is accepted and consumed and, in backward paddle unit motion, the paddle's energy is accepted and consumed. This is the second technological innovation of my patent application over patents: GB 384603A and JP 57081168.    5. U.S. Pat. No. 5,244,359 describes a marine energy converter system based on utilization of differences in sea level which it converts into linear movement only of piston 21 which moves within its casing 17. The casing 17 is rigidly connected with pole 3 which is anchored at the bottom of the sea. The above converter can only exploit linear monement, and does not make it possible to exploit the dynamic energy of the sea waves. The above converter is equipped with a floating system piston. The piston is connected to a float 41 and this enables the linear movement of the piston within the piston casing. As a result the piston, within its casing, by virtue of the rising level of water, requires very little force to rise. However, since energy is always a product of force multiplied by distance, the resultant energy is very small as a consequence of the small amount of force exerted.\nThe present invention, however, is based on a different principle, i.e., the exploitation of the dynamic energy of sea waves, which give rise to pivotal movement of the paddles rather than their elevation. The pivotal movement is subsequently converted into linear movement of the piston by means of a crankshaft mechanism 32,40. The piston rises within the casing with relatively great force and the energy produced is thus greater.\nThanks to the structure of the system according to the present invention, it is possible to utilize waves regardless of the height of the wave itself.\nBoth from the point of view of the present invention's operating principle and from the point of view of its particular structure, the U.S. Pat. No. 5,244,359 energy converter is not similar to the present invention.    6. U.S. Pat. No. 4,843,249 describes a hydroelectric system, which utilizes the movement of the waves and converts them into circular movement. Said system includes a turbine 22 whose vanes are moved by the moving wave. The turbine wheel is connected to equipment for the generation of electrical energy.\nThe present invention, by comparison, does not convey circular energy, but rather utilizes the axial movement of the paddles to move the piston in a linear motion.\nAs in the previous case, the above U.S. patent does not void the innovation or the innovation of the present invention, neither in principle nor from the point of view of the structure of the above-described patent.\nIt may also be noted, that the utilization of wave energy in order to create circular energy along, is inferior to the creation of linear motion, on the following counts:\nIn order to achieve better that 90% utilization, for example, a pump activated through circular energy must rotate as speeds of 500-5000 rpm. Below 500 rpm, the efficiency rate is abruptly reduced to about 40%. In order to attain such high speeds, when the wave movement is relatively low, it is necessary to provide additional transmission by means of a few pairs of cog-wheels, thus creating losses in the system. Moreover, if we introduce a number of pairs of cog-wheels, then the total efficiency of the system is reduced, and the energy produced is thus far lower. The cost of the system too, is increased through the provision of may pairs of additional cogwheels.\nHowever, when the piston moves in a linear motion, the immediate efficiency of the system is at least 90% and a linear pump is not limited by any constraints.    7. U.S. Pat. No. 53,311,064 describes a system for generating energy from the movement of sea waves, based on a similar principle to that of the above described hydro-electric system, i.e., the use of a turbine which is rotated by the flow of water. The circular movement is conveyed to the energy generating mechanism by means of a transmission mechanism.\nAll the above given explanations regarding the hydro-electric system according to U.S. Pat. No. 4,843,249, are also valid in this case, both as regards circular motion versus linear motion and the differing structure.    8. Fr Patent 501795 described a system for generating energy from the movement of sea waves.\n(i) The system has one paddle. The low extremity of the paddle's plate is a buoy that moves vertically up and down when the wave reaches the plate. The paddle receives only the height energy from the wave.\nWhen a wave hits a paddle, the paddle goes up in pivotal movement. In this way, each paddle receives energies from: a. the wave height energy which is, potential energy plus b. the wave movement energy, which is kinetic energy. The present invention receives from each paddle much more energy in comparison to the above Fr Patent 501795.\n(ii) The present invention has at least two paddles. The Fr Patent 501795 has one paddle. The advantage of the present invention from this reason is the same as explained in 2(i) above as regards towards U.S. Pat. No. 5,084,630.    9. JP patent 006750 converters can only exploit linear movement, and does not make it possible to exploit the dynamic energy of the sea waves. The hydraulic cylinders 2 and 2a, are pistons connected to a float 1 and this enables the linear movement of the pistons within the piston casing. As a result the pistons within the casing, by virtue of the rising level of water, require very little force to rice. However, since energy is always a product of force multiplied by distance, the resultant energy is very small as a consequence of the small amount of force exerted.\nThe present invention, however, is based on a different principle, i.e., the exploitation of the dynamic energy of sea waves, which give rise to pivotal movement of the paddles rather than their elevation. The pivotal movement is subsequently converted into linear movement of the bi-directional piston by means of a crankshaft mechanism 32, 40. The piston rises within the casing with relatively great force and while the energy produced is a product of force multiplied by distance, the resultant energy is higher as a consequence of the high amount of force exerted. Thus the energy produced according to this system is greater then JP patent 006750.    10. U.S. Pat. No. 4,490,621 described a system for generating energy with two hydraulic motors, one hydraulic motor 121a, for low waves and one hydraulic motor 121b for high waves with very complicated system. When the oil pressure PLa>PLb then first hydraulic motor 121a works and when the oil pressure PLa<PLb, the second hydraulic motor works. The present invention works with all kinds of oil pressure with the same hydraulic motor. If for example, the present invention works with a system of three hydraulic motors, the U.S. Pat. No. 4,490,621 requires 6 hydraulic motors: three for law waves plus three for high waves, which is very expensive and with very complicated system. The present invention is more economical and much simpler than U.S. Pat. No. 4,490,621.    11. U.S. Pat. No. 1,960,622 generate energy does not establish optimum distance of the paddle heights, to make plan maximum energy production, since the rod 20 and the plate 25 move upwards or downwards synchronously, automatically, but not supervised as described in the present invention, with the movement of the sea level at tide time.            The paddles' system according to the present invention stays in the optimum fixed position even when the sea water continues to move upwards and downwards with the tide to create the plant maximum energy production, while in the U.S. Pat. No. 1,960,622, the rod 20 and the plate 25, move upwards or downwards all the time with the sea water at the tide time.        "}
{"text": "The present invention relates to a magnetic disc device used in computers, information storage devices and so on, a magnetic storage device used in such information home appliances as digital VTRs, and a magnetic recording, and and, more particularly, to a magnetic recording and reading device suitable for realizing high-speed recording and reading, and for high-density recording.\nSemiconductor memories, magnetic memories, etc., are used in the storage or recording devices of information equipment. Semiconductor memories are used in internal primary storage in the light of high-speed accessibility and magnetic memories are used in external secondary storages in the light of a high capacity, low cost and nonvolatile property. Magnetic disk devices, magnetic tapes and magnetic cards are the main current in the magnetic memories. A magnetic recording portion which produces a strong magnetic field is used in order for writing magnetic disks, magnetic tapes or magnetic cards. Further, reading portions based on the magnetoresistance effect or the electromagnetic induction effect are used in reading magnetic information recorded at a high density. In recent years, for reading, the giant magnetoresistance effect and the tunneling magnetoresistive effect have also begun to be examined. These functional portions for recording and reading are both installed in an input-output part which is called a magnetic head.\nThe basic configuration of a magnetic disk device is shown in FIGS. 10A and 10B. FIG. 10A shows a plan view of the device and FIG. 10B shows a vertical-sectional view of the device. Recording media 101-1 to 101-4 are fixed to a hub 104 to be rotated by a motor 100. In FIG. 10B shows one example which comprises four magnetic disks 101-1 to 101-4 and eight magnetic heads 102-1 to 102-8. However, the magnetic disk device may comprise at least one magnetic disk and at least one magnetic head. The magnetic heads 102-1 to 102-8 move on the rotating recording media. The magnetic heads 102-1 to 102-8 are supported by a rotary actuator 103 via arms 105-1 to 105-8. Suspensions 106-1 to 106-8 have function of the pressing the magnetic heads 102 against the recording media 101-1 to 101-4 under a determined load, respectively. A given electric circuit is needed for processing of reproduction signals and for inputting and outputting of information. Recently, a signal processing circuit in which waveform interference at high-density is positively utilized, such as PRML (Partial Response Maximum Likelihood) or EPRML (Extended PRML) which is an enhanced. PRML, has been adopted, contributing greatly to a high-density design. The signal processing circuit is installed in a circuit board on a cover 108, etc.\nThe functional portion for writing and reading information on a magnetic head assembly is comprises components shown in FIG. 11A, for example. A writing portion 111 is comprised of a spiral coil 116 between magnetic poles 117, 118 which are magnetically connected with each other. The magnetic poles 117, 118 are both composed of a magnetic film pattern, which are made of an NiFe alloy, etc., respectively. The reading portion 112 comprises a magnetoresistance element 113 made of an NiFe alloy, etc. and an electrode 119 for applying a constant current or a constant voltage to the element 113 and for detecting changes in resistance. The magnetic pole 118, which is made of an NiFe alloy, etc. and serves also as a magnetic shielding layer, is provided between the writing and reading portions. There is further a shielding layer 115 underneath the magnetoresistance element 113. A reading resolution is determined by the clearance distance between the shielding layer 115 and the magnetic pole 118 (serving also as another shielding layer). The functional portion is formed on a magnetic head slider 1110 (FIG. 11B) via an underlayer 114 made of Al2O3, etc. Incidentally, the magnetic head slider, which is provided with a protection layer made of hard-carbon, etc. on the surface opposed to the magnetic recording medium, is supported by a gimbal 1111 and a suspension 1113, as shown in FIG. 11B. The magnetic head slider moves relatively to the magnetic recording medium while floating from the medium surface and, after positioning in an arbitrary position by an arm 1114 connected to a motor, realizes the function of writing or reading magnetic information via lead lines 1116 and 1115. With respect to the above function, there is also provided an electric control circuit together with the aforementioned signal processing unit or on the head carriage.\nA detailed structure of a recording medium is schematically shown in FIG. 12. As described in JP-A-3-16013, most of the conventionally used recording media are produced by forming a magnetic layer 123 made of a Co—Cr—Ta alloy, or a Co—Cr—Pt alloy, etc. on a non-magnetic substrate made of Al plated with an NiP alloy, a glass, a high-hardness ceramics, a polished Si or the like, or a plastic substrate 121 by the sputtering method, or the evaporation method, or the plating method, etc. Usually, an under layer 122 made of Cr, or a Cr alloy, etc. for orientation control of the magnetic layer is often formed on the substrate. Furthermore, a protection film 124 made of diamond-like carbon containing nitrogen and/or hydrogen, or SiO2 or SiN or ZrO2, etc. is provided to ensure durability of sliding resistance, and a lubricating film 125 made of perfluoro alkyl polyether having an adsorptive or a reactive end group, or organic fatty acids, etc. is provided.\nIn addition to the magnetic recording device, magneto-optic recording devices that perform recording and reading on a magnetic recording medium through the use of light have also been put to practical use. The magneto-optic recording devices are classified into one type in which recording is performed only by light modulation and another type in which recording and reproduction are performed by light with a modulated magnetic field. However, the both types greatly rely on heat when recording and reading. Therefore, according to such type of devices, it is impossible to perform recording and reading in high data transfer rate and thus they have been adopted mainly in backup systems, etc.\nThe importance of a storage device is determined by its storage capacity and the speed during inputting-outputting operations. In order to increase competitiveness of products, it is necessary for the storage device to increase capacity by higher recording density, higher rotational speed and higher data transfer rate than those of the prior art. Thus, an important problem to be solved by the present invention is to provide a device capable of recording and reading at a high data transfer rate of not less than 50 MB/s and, more preferably, that at a high density of not less than 5 Gb/in2. A magnetic recording medium capable of recording and reading at a high frequency and capable of obtaining a high S/N ratio at a high density and a magnetic head capable of generating a sufficient magnetic recording field at a high frequency are necessary for meeting the requirement.\nIn conventional magnetic recording media, there have been proposed and actually carried out to reduce noise by refining crystal grains in order to obtain a high S/N ratio at a high density of about 1 to 3 Gb/in2, and by promoting segregation of non-magnetic components at grain boundaries to reduce exchange coupling among crystal grains as being taught in JP-A-63-148411, JP-A-3-16013 and JP-A-63-234407 so as to make the coercive squareness S* to not more than 0.85 and the rotational hysteresis loss RH to the range of 0.4 to 1.3. Noise can be considerably reduced by recording and reading at a data transfer rate of not more than about 20 MB/s. However, when the magnetic recording was carried out on that film media of the prior art at a high frequency of not less than 50 MB/s, thermal fluctuation effects in fine magnetic crystallines is remarkable due to weak exchange coupling among crystal grains and the apparent coercive force is high resulting in that it was impossible to record on it accurately. Furthermore, even when recording is performed under a large current with utilization of a modified recording circuit, etc., the magnetic recording transition region is widened due to a broad magnetic recording field resulting in that noise increases and/or recorded information is lost when it was alowed to stand for a long time."}
{"text": "The combination of solid-state light sources like light-emitting diodes (LEDs) and lasers, with microdisplays like liquid-crystal-on-silicon (LCOS) devices is leading to the development of electronic projectors small enough to be embedded in, for example, the handset of a mobile phone. In such “pico projectors,” small optical engine size and high optical efficiency are especially important. Achieving size and efficiency goals requires new ways of managing polarized light to illuminate and image microdisplay panels. It would increase the efficiency of such optical systems if they could efficiently use unpolarized light to illuminate a display that ordinarily requires polarized light, particularly if this could be accomplished without doubling the étendue of the light source.\nIt is against this background that the teachings disclosed herein have been developed."}
{"text": "FIG. 1 illustrates a prior art HDD drive amplification writer circuit 100. Generally, the circuit 100 sources large switching currents (from “Iw/X”, 110, 112, 150, 152, wherein “X” is ratio of an input resistor 120, 122 to a sum of output resistors 130, 132) through a resistive load 120, 122 to create the Iw reference voltage into a corresponding impedance match buffer amplifier 115, 117. In the circuit 100, current is sourced or disabled through the various current sources Iw/X, 110, 112, 150, and 152, depending upon whether it is desired to write a “1” or a “0” to a hard drive. In other words, a current can be driven through input resistors 120, 122, in opposite directions, and a voltage associated with the input resistor 120 at DC_P, and the input resistor 122 at DC_N drives inputs of the amplifiers 115, 117.\nIn the circuit 100, current sources 112, 150, are labeled PHI1 and current sources 110, 152 are labeled PHI2. In the circuit 100, PHI1 switching current 112 and PHI1 switching current 150 would be enabled and PHI2 switching current 110 and PHI2 switching current 152 would be disabled, or alternatively, PHI1 switching current 112 and PHI1 switching current 150 would be disabled and PHI2 switching current 110 and PHI2 switching current 152 would be enabled, whichever was preferred to write a “1” or a “0” between the write nodes 135 and 137.\nHowever, there is a fundamental tradeoff between characteristics of a switching current, provided by current sources Iw/X 110, 112, 150, 152 and a switching speed (i.e. how quickly a transition can occur between a “1” or a “0”) of the preamplifier HDD write circuit 100. In order to achieve a desired switching speed, the resistance of the input resistors 120, 122 was kept small, such as a ratio of four times the sum of the output resistors 130, 132 and the load resistance to drive the output write nodes 135, 137. Consequently, with a gain of only four times to an output write nodes 135, 137 the switching current Iw/X of the current sources 110, 112, 150, 152, then was consequently disadvantageously large to achieve the needed write current magnitude at the output write nodes 135, 137. This is a high power consumption approach. For example, for a 60 milliamp write current, a constant 15 milliwatts is being variously sourced by sources 110, 112, 150, 152. Moreover, these current sources would be switched on or off, depending upon the value desired to be written to the hard drive, which further consumed power and created transients.\nFurthermore, even with employing a smaller value of a current Iw/X so as to decrease power consumption, this introduces yet other disadvantages. The resistors 120, 122 need to be larger with smaller currents Iw/X to achieve the same needed write current magnitude at the write output nodes 135, 137. The input resistors 120, 122 create an additional RC pole, along with the parasitic capacitance at a node of DC_P and DC_N, and the capacitance can be large due to large devices and large metal capacitance to support the large switching current. As the input resistors 120, 122, become larger, this RC pole becomes more dominant within the circuit 100 and slows the switching speed.\nYet still further, the switching current is “X” dependent, as the switching current is employed to be gained up/amplified by a ratio of output resistor 130 divided by input resistor 120 and output resistor 132 divided by input resistor 122 to create the output write current. These ratios not only force the switching currents to be large in order to generate a large output write current as discussed above, but also does not allow for slew rate control at the write nodes DC_P, DC_N and hence at write nodes 135, 137. It would be desirable to be able to vary the slew rate at the write nodes 135, 137. However, doing so in this architecture would also vary the output write current, which needs to remain at the same value for all slew rates.\nTherefore, there is a need in the art to address at least some of the issues associated with conventional HDD controllers."}
{"text": "Known porous bodies employed for filter materials or catalytic carriers include those consisting of various materials such as resin, metals or ceramics. Among these, a filter or a catalytic carrier consisting of a ceramics material is generally employed in high temperature or strongly corrosive environments which other materials cannot withstand. A filter or a catalytic carrier consisting of oxide ceramics such as alumina (Al.sub.2 O.sub.3) has already been put into practice.\nAs to a porous body consisting of nonoxide ceramics, on the other hand, only small examples have been put into practice while Japanese Patent Laying-Open No. 63-291882 discloses a silicon nitride based or silicon carbide based porous body prepared by a heat treatment. Further, Japanese Patent Laying-Open No. 1-188479 discloses a method of compacting mixed powder of silicon powder and silicon nitride powder of relatively coarse particles and thereafter nitriding the same thereby preparing a porous body as a solid target.\nAs hereinabove described, it is difficult to use a porous body consisting of resin or a metal in a high temperature or corrosive atmosphere. It is inevitably necessary to employ a porous body made of ceramics for a filter for removing foreign matter from a high-temperature exhaust gas or for a carrier serving as a catalyst for decomposing a harmful matter.\nAs an example of such porous bodies made of ceramics, porous bodies made of alumina have been put into practice. While the porous bodies of alumina are varied in pore size, porosity and bending strength, a porous body having a porosity of 35 to 40% and a mean pore size of 25 to 130 .mu.m has a bending strength of 20 to 35 MPa, whereby the strength of the porous body is insufficient depending on its use.\nIn the silicon nitride based porous body disclosed in the aforementioned Japanese Patent Laying-Open No. 63-291882, the porosity is less than 30% and fluid permeability is insufficient. In general, the strength of ceramics tends to be reduced following an increase in the porosity, and it has been extremely difficult to attain compatibility between the porosity and the strength."}
{"text": "The present invention relates to methods for detecting and measuring plasmacytoma cells in a sample.\nPlasmacytoma, represented by multiple myeloma, is generally termed a myeloma, and is pathologically divided into multiple myeloma, plasma cell leukemia, solitary plasmacytoma, extramedullary plasmacytoma, smoldering multiple myeloma, or asymptomatic myeloma. All of them are B cell tumors and are characterized by abnormal growth of monoclonal plasma cells (antibody-producing cells). Many of these tumors accumulate in the bone marrow and are diseases of poor prognosis accompanied by systemic bone lesions.\nIn multiple myeloma, terminally differentiated B cells or plasma cells that produce and secrete immunoglobulins monoclonally increase predominantly in the bone marrow, and thereby monoclonal immunoglobulins (M proteins) or their constituent light chains and/or H chains are detected in the blood and urine. With the aging of the population in recent years, patients with multiple myeloma are growing in number, and the mortality rate has reached 2.1 in 100,000 people in Japan.\nThe diagnosis of multiple myeloma has been made based on the detection of M proteins and/or bone lesions. These methods, however, required that the number of myeloma cells to be detected amounted to 0.2 to 1.2xc3x971012, systemically, and even for the detection of M proteins, growth of myeloma cells to about 1.2xc3x971010 was required. This rendered the early detection thereof difficult.\nFurthermore, since these methods of diagnosis do not always reflect the amount of tumor cells, it was inconveniently impossible to decide effective regimens for treatment. In addition, it was also difficult to differentiate from benign monoclonal immunoglobulinemia in which M proteins are similarly detected, or to diagnose the non-secretory type myeloma in which no M proteins are present.\nOn the other hand, there have been attempts to diagnose myeloma using antibodies that recognize antigen on the cell membrane of plasma cells. These antibodies include such monoclonal antibodies as anti-PCA-1 (Anderson, K. C. et al., J. Immunol. (1983) 130, 1132), anti-PC-1 (Anderson, K. C. et al., J. Immunol. (1983) 132, 3172), anti-MM4 (Tong, A. W. et al., Blood (1987) 69, 238), anti-CD38 (Epstein, J. et al., N. Eng. J. Med. (1990) 322, 664, Terstappen, L. W. M. M. et al., Blood (1990) 76, 1739, Leo, R. et al., Ann. Hematol. (1992) 64, 132, Shimazaki, C. et al., Am. J. Hematol. (1992) 39, 159, Hata, H. et al., Blood (1993) 81, 3357, Harada, H. et al., Blood (1993) 81, 2658, Billadedeau, D. et al., J. Exp. Med. (1993) 178, 1023) and the like. Since such detection using antibody required a sample that contained a large amount of myeloma cells, it was particularly unsuitable for early detection in which myeloma cells were not always present.\nOn the other hand, Goto, T. et al. reported that they immunized mice with human plasma cells and obtained a mouse monoclonal anti-HM1.24 antibody that recognizes an antigen having a molecular weight of 29-33 kDa specifically expressed on B cell lines (Blood (1994) 84, 1922). It was further reported that the antigen (HM1.24) recognized by monoclonal anti-HM1.24 antibody is considered to be an antigen associated with the terminal differentiation of B cells (Goto, T. et al., Jpn. J. Clin. Immun. (1992) 16, 688), and that monoclonal anti-HM1.24 antibody reacts to plasmacytoma in a specific manner (Shuji Ozaki et al., the program of The 19th General Meeting of Japan Myeloma Study Group, General presentation 3). It was not known, however, that the diagnosis of plasmacytoma could be effected using the gene of HM1.24.\nIt is an object of the present invention to provide a method that permits a simple and sensitive detection or measurement even with a small amount of tumor cells, and that permits the diagnosis of plasmacytoma.\nIn order to attain the above object, the present inventors have devised a combination of primers and a probe based on the already-known sequence of the HM1.24 gene (Japanese Patent Application No. 9-271536), attempted to detect and measure the amount expressed of HM1.24 mRNA in a simple and sensitive manner and to quantitate the amount expressed. As a result of intensive research, the present inventors have found that plasmacytoma cells can be specifically detected or measured by PCR-amplifying HM1.24 mRNA from a small amount of tumor cells and using a probe labeled with fluorescence, and that plasmacytoma can be diagnosed by amplifying similarly normal cells or tumor cells other than plasmacytoma cells as a control sample and then comparing the amount of the amplified product, and thereby have completed the present invention.\nThus, the present invention provides a method of detecting or determining plasmacytoma cells in a sample, said method comprising amplifying a polynucleotide that is specifically or strongly expressed in plasmacytoma cells and then detecting or measuring the amplified product.\nAs said plasmacytoma, there can be mentioned multiple myeloma, plasma cell leukemia, solitary plasmacytoma, extramedullary plasmacytoma, smoldering multiple myeloma, asymptomatic myeloma and the like.\nPreferably, said polynucleotide is DNA or mRNA, in particular mRNA, that encodes the HM1.24 antigen.\nPreferably, said amplification method is the PCR method.\nPreferably, said measurement method is a method that employs a labeled probe.\nPreferably, said label is a fluorophore, a radioisotope, an enzyme or a combination thereof.\nThe present invention also provides a method of diagnosing plasmacytoma by detecting or measuring plasmacytoma cells in a sample, said method comprising amplifying a polynucleotide that is specifically or strongly expressed in said plasmacytoma cells, detecting or measuring the amplified product, and then comparing the amount of the amplified product to that in a control sample.\nPlasmacytoma that is the subject of the detection or measurement method of the present invention is a tumor that expresses the HM1.24 antigen having the amino acid sequence as set forth in SEQ ID NO: 1 or the amino acid sequence substantially identical to said amino acid sequence, and includes multiple myeloma, plasma cell leukemia, solitary plasmacytoma, extramedullary plasmacytoma, smoldering multiple myeloma, asymptomatic myeloma and the like.\nSamples containing plasmacytoma cells may be prepared by isolating leukocytes from the peripheral blood or the bone marrow of patients with plasmacytoma and collecting RNA therefrom. Since many of the myelomas tend to form a tumor mass in the bone marrow, it is desirable to collect leukocytes containing plasmacytoma cells by the bone marrow puncture to said tumor mass. In cases where there is a marked appearance of tumor cells in the peripheral blood as in plasma cell leukemia etc., the use of peripheral blood is preferred because of simplicity in handling.\nIn order to isolate leukocytes containing plasmacytoma cells from the peripheral blood or the bone marrow collected from patients with plasmacytoma, a commercially available reagent may be used for separation by gradient centrifugation. As the commercially available reagent, there can be used Ficoll-Paque (manufactured by PHARMACIA BIOTECH), a Mono-Poly Resolving Medium (manufactured by Dainippon Pharmaceuticals), and the like. It is also possible to prepare leukocytes by lysing red blood cells in a hypotonic buffer and then washing by centrifugation. As a solution used in washing, there can be used isotonic solutions such as PBS and MEM.\nIn order to prepare RNA from the thus obtained leukocytes, a method well known to those skilled in the art may be used (see, for example, Idensi Sousa Jikkenhou (Experiments on Genetic Manipulation), Yasutaka Takagi, Kodansha Scientific, pp. 24-39, 1980). RNA may be total RNA or mRNA. Total RNA or mRNA can be extracted using a commercially available kit, for example TRIZOL (manufactured by GIBCO BRL), ISOGEN (manufactured by Wako Pure Chemicals Industries, Ltd.), and mRNA Purification Kit (manufactured by PHARMACIA BIOTECH).\nAny target polynucleotide can be used as long as it is specifically or strongly expressed in plasmacytoma cells. However, due to ease of detection, it is preferably DNA or mRNA, in particular mRNA, encoding the HM1.24 antigen that is specifically expressed in plasmacytoma cells.\nIn accordance with the present invention, a specific example of a strongly expressed polynucleotide includes a polynucleotide that is expressed five-fold or more, preferably 10-fold or more, and most preferably 20-fold or more than normal cells or tumor cells other than plasmauytoma cells.\nA relative amount of the polynucleotide expressed in a sample may be determined by subjecting RNA or mRNA extracted from said cells to agarose gel electrophoresis, allowing the product to be adsorbed to a membrane, performing Northern blot hybridization using a radio-labeled probe, and then determining and comparing the radioactivity. Alternatively, the competitive RT-PCR method (Competitive RNA Transcription Kit, Takara) may be used to amplify, at the same time, the RNA competitor and the mRNA of interest in the same reaction mixture with the same primer. Then, based on the ratio of the amplified products, the original amount of mRNA can be estimated. Furthermore, it can be simply determined using ABI PRISM7700 according to the method described in Examples.\nAs the method of amplification, the PCR method (polymerase chain reaction method) can be used. When the polynucleotide of interest is mRNA, the use of RT-PCR is preferred. For the PCR method or the RT-PCR method, a commercially available reagent may be used including the TaKaRa RNA LA PCR(trademark) Kit (AMV) Ver 1.1 (manufactured by Takara Shuzo), the RT-PCR high-Plus-(manufactured by Toyoboseki), or the TaqMan EZ RT-PCR Kit (manufactured by Perkin-Elmer) and the like.\nAs a method of amplification, in addition to the PCR method, the TMA method (Transcription Mediated Amplification method: Japanese PCT Publication No. 4-500759) that amplifies RNA may be used.\nThe primer is a pair of oligonucleotides that can amplify a segment of the nucleotide sequence as set forth in SEQ ID NO: 1, specifically a pair of oligonucleotides that can hybridize, at a certain interval, to a nucleic acid having the nucleotide sequence as set forth in SEQ ID NO: 1 or a nucleic acid complementary thereto. The distance on the sequence of SEQ ID NO: 1 to which the pair of oligonucleotides hybridize is, but is not limited to, a range of 30 to 400 nucleotides, preferably 50 to 200 nucleotides.\nThe primer is preferably selected from the region in which the amino acid sequence is well conserved among the species and mainly in the extracellular region from positions 48 to 124 in the amino acid sequence of SEQ ID NO: 1.\nBased on the above conditions, the primer can be selected using, for example, a commercially available analytic al software such as Primer Express (manufactured by Perkin-Elmer).\nThe length of the primer to be used for amplification is, but is not limited to, 10 to 100 nucleotides, preferably 10 to 50 nucleotides, more preferably 15 to 30 nucleotides, and more preferably 20 to 25 nucleotides, for example about 20 nucleotides.\nSpecific examples of such a primer pair include the forward primer (SEQ ID NO: 3) (corresponding to positions 153 to 172 in the nucleotide sequence as set forth in SEQ ID NO: 1) and the reverse primer (SEQ ID NO: 4) (corresponding to positions 304 to 323 in the nucleotide sequence as set forth in SEQ ID NO: 1).\nThe probe is an oligonucleotide that specifically hybridizes to a nucleotide sequence defined by said one pair of primers. Thus, such a probe is selected from the region between the forward primer and the reverse primer. Though the positions on SEQ ID NO: 1 corresponding to the probe depend on the sites to which the primer pair hybridizes, those that hybridize to the region of nucleotide number 173 to 303 in SEQ ID NO: 1, preferably the region of nucleotide number 224 to 249 are preferred. The length of the probe to be used is, but is not limited to, 10 to 200 nucleotides, preferably 20 to 100 nucleotides, for example 25 to 50 nucleotides, for example about 30 nucleotides.\nThe amplified polynucleotide can be detected or measured by subjecting the PCR reaction mixture to agarose gel electrophoresis in the presence of a luminescent substance such as ethidium bromide and then measuring the light emitted by said band. When a probe is used, the polynucleotide is detected by allowing it to be adsorbed to a nitrocellulose membrane followed by hybridization to the labeled probe. The sequence or the length of the probe is not limited as long as it can specifically detect the target polynucleotide to be detected.\nFor labeling, a fluorogenic substance, a radioisotope, an enzyme or a combination thereof may be used. As the fluorogenic substance, FAM (6-carboxyfluorescein), JOE (6-carboxy-4,5-dichloro-2,7-dimethoxyfluorescein), TET (6-carboxy-4,7,2xe2x80x2,7xe2x80x2-tetrachlorofluorescein), and HEX (6-carboxy-4,7,2xe2x80x2,4xe2x80x2,5xe2x80x2,7xe2x80x2-hexachlorofluorescein) may be used. As the radioisotope, 32P may be used.\nAs the enzyme, alkaline phosphatase may be used, and as the substrate for the enzyme, CDP-Star(trademark) (manufactured by BOEHRINGER MANNHEIM), for example, may be used. By using ABI PRISM7700 developed by Perkin-Elmer that permits the determination of the product with time while performing the PCR reaction, simple determination can be attained. As the reagent for RT-PCR in this case, TaqMan EZ RT-PCR Kit (manufactured by Perkin-Elmer) is preferably used."}
{"text": "The present invention relates to retail store checkout stands, and more specifically, to checkout stands having an auxiliary structure that allows face to face contact between the cashier and the customer and that also allows vertical adjustment of the structure.\nAnyone who has been in a retail store is familiar with the checkout stand; the checkout stand is where one pays for the selected items. Checkout stands, also commonly referred to as checkstands, can be arranged in a broad variety of configurations. The checkout stand may include one belt, two belts, three belts, or merely have a stationary surface, generally supported by a counter or cabinet. A bar code scanner is typically recessed into the counter or cabinet. Also included at the checkout stand are the register, cash drawer, a keyboard, a credit card machine (often referred to as a credit card swipe and pin pad machine), a receipt printer, monitor or display, telephone, and other such accessory equipment. The register, cash drawer, and other equipment is typically positioned on a cabinet often referred to as a cashstand, which is positioned to the left or right of the cashier, at a right angle, as the cashier is looking at the customer.\nThe goal within the industry has been to arrange the checkout stand in a manner that is convenient for both the customer and the cashier. The layout of the checkout stand and the configuration of the equipment impact both the customer and cashier. Because the checkout stand is a highly visible feature of the retail shopping experience for the customer, providing a well arranged and configured checkout stand is important. The layout of the checkout stand can enhance or deter from the customer\"\"s experience, which can affect whether or not that customer shops at that store again. For example, a cluttered checkout stand, with cables and electrical cords exposed, is messy and uninviting. A checkout stand not well organized is also viewed as messy and as inefficient.\nFor the cashier, standing at the checkout stand is a daily occurrence, with repetitive motions. When the cash drawer and register are located at a cashstand, the cashier has to twist or rotate 90 degrees from the scanner to the register and back to the customer, losing both eye contact with the customer as well as their view of the cash drawer, which may be inadvertently left open as they again turn to face the customer.\nWith more and more equipment being added to checkout stands, the space available for positioning the equipment is at a premium. One solution has been to include a podium, console, or other auxiliary structure positioned away from the cashstand, typically positioned between the cashier and the customer in the vicinity of the bar code scanner. In some designs, these consoles are built onto a counter surface; in other designs, the consoles are bolted or otherwise attached to the counter. This console generally supports some of the equipment integral to the checkout process, for example, such as the keyboard, a credit card swipe and pin pad machine, and/or a display monitor. A separately mounted horizontal pad, often called a checkwrite, can be used by the customer for writing checks or for signing credit card receipts. The cash drawer and receipt printer have, to date, remained to the left or right of the cashier at the cashstand.\nThe arrangement of having the console above the scanning area and above the scanner is generally a superior arrangement. To maximize the ease of the checkout transaction, it is desired to keep the height of the console as low as possible, preferably with near zero clearance between the tip of the scanner and the bottom of the console. However, lowering the height of the console to improve customer-cashier view can hinder the minimum clearance needed above the scanner, for example, for removal of the scanner for maintenance and service.\nWhat is desired is a compact console arrangement that allows access to available equipment, such as the scanner, and provides a configuration that is friendly and convenient to both the customer and cashier. What is also desired is a console arrangement that minimizes cashier fatigue and injury potential due to twisting stresses, that provides cash drawer security, and that speeds the checkout transaction.\nThe present invention is directed to a console for use with a store checkout stand. The console of the present invention reduces twisting and turning of the cashier by providing the cash drawer in a location between the cashier and the customer. The console also incorporates an elevation system to facilitate raising and lowering the console top, on demand, for example, to gain access to equipment, such as the bar code scanner, positioned below the console top.\nThe console can be incorporated into any type of checkout stands, such as with conveyor belt systems or with a checkout stand that simply uses a counter. The console can be mounted onto a counter or cabinet, such as on a siderail or sidewall, on a vertical or a horizontal surface.\nIn one particular aspect, the console comprises a stand support structure and a top support structure, the stand support structure having a leg, a mounting portion fixed to and movable in respect to the leg, and a biasing element connected to the leg and to the mounting portion. The top support structure is attached to the stand support structure and has a first portion, and a second portion, the second portion movably attached to the first portion. The mounting portion of the stand support structure can be a bracket, configured for mounting to a vertical surface such as a counter sidewall. Alternately, the bracket can be configured for mounting to a horizontal surface, such as a floor. The biasing element can be a spring, such as a coil or coiled spring. The first portion can have a top fixedly attached thereto, and the second portion, which is movable in respect to the first portion, is fixedly attached to a surface, such as a vertical counter sidewall. The first portion of the top support structure can have at least one slot therein, and the second portion has an element for engaging the slot. In most embodiments, the first portion will have a slot at each end. In an alternate embodiment, the second portion can have at least one slot therein, and the first portion has an element for engaging the slot.\nIn another particular aspect, the leg, which, by means of a saddle, attaches to the first portion of the top support structure. It has been found that the weight of the console plus the cash drawer, its contents, and any other equipment, create considerable torque among the top support structure and the stand support structure. The saddle is constructed to transfer force from the console top to the leg and mounting portion.\nPreferably, the console is mounted between the cashier and the customer, above or in close proximity to the bar code scanner. In a preferred embodiment, the console includes a volume for receiving a cash drawer. By positioning the cash drawer in a location in front of the cashier when the cashier is facing the customer, twisting and turning of the cashier is minimized, and the drawer is placed in front of the cashier, rather than at a side location where security can be an issue. Further, by positioning the cash drawer in front of the cashier, the efficiency of the checkout transaction is increased.\nOther console constructions, checkout stands incorporating the console, and methods of raising and lowering the top of a console, are disclosed."}
{"text": "1. Field of the Invention\nThe present invention generally relates to current limiting contact systems which are typically used in circuit breakers. In particular, the present invention relates to a contact system, which contact system is composed of several parallel contact fingers for each pole to be switched and each pole to be protected, and, when there is a short circuit, limits the short circuit current by means of electrodynamic opening, until the breaker mechanism is tripped.\n2. Background Information\nA three-phase circuit breaker with a contact system of the type described above is disclosed in European Patent Application No. 353 948 A2. Three contact systems are arranged next to one another in a molded housing. Each contact system is composed of a stationary contact with a stationary contact piece, and a moveable contact with moveable contact pieces. Each stationary contact leads into a first phase connector. The moveable contact is composed of several contact fingers arranged parallel to one another, each with a moveable contact piece. The contact fingers are mounted on an axle, and each contact finger is connected by means of a V-shaped, flexible connection lead to the connection lead. The connection leads are connected to a second phase connector.\nEach contact finger is acted upon by a helical compression spring, which compression spring is braced against a contact carrier. The contact carrier grips longitudinally in a U-shape around the contact fingers, and the axle of the contact fingers is mounted in the contact carrier. The contact fingers are mounted on a transverse selector shaft by means of a swivel pin. Each contact carrier is acted upon by a cam-operated toggle mechanism. In addition, two helical tension springs serve as energy stores or energy storage mechanisms, and engage between the swivel pin and the axles of cam follower rollers, which cam follower rollers are mounted in a linked manner in the switch axle. The central moveable contact is additionally connected to the toggle system of a breaker mechanism via the selector shaft by means of lateral supports. The necessary contact pressure force between the moveable contact pieces and the stationary contact piece is produced by the helical compression springs. During tripping or manual shutting off by means of the breaker mechanism, essentially all moveable contacts are swivelled by their spring-activated coupling with the toggle spring mechanism over the selector shaft, in order to open the contact pieces.\nIn case of extreme short circuiting, the electrodynamic repulsive forces caused by the short circuit current flowing in the V-shaped connecting leads, as well as the high current density on the contact pieces, exceed the contact compression force which leads to spinning off of the moveable contact concerned. The toggle spring mechanism then falls into its second stable tilted position. This tilted position is maintained until the circuit breaker is opened by the subsequent tripping breaker mechanism in all phases. The arc occurring between the separated contact pieces is extinguished with additional arcing contacts and an arcing chamber. Elastic stops are housed in the molded housing, which on the one hand collect the high kinetic energy of the electrodynamically spun-off contact concerned, and on the other hand tilt this contact back into the first stable position after tripping of the breaker mechanism.\nThe disadvantages of the contact system described above are, on the one hand, the complex means for the conductive connection between the contact fingers and the second phase connector the complex means; for the production of the necessary contact compression force the complex means; and for the collection of the kinetic energy of the spun-off contact; and, on the other hand, the considerable number of elements which are necessary for transferring movement from the breaker mechanism to the moveable contact. A further, crucial disadvantage is that with increased wear of the contacts, in particular by burning, the contact compression forces diminish, leading to an increasing degradation of the switching properties.\nIn German Patent Application No. 31 33 285 C2 an automatic circuit breaker is disclosed. This automatic circuit breaker includes a moveable contact carrier mounted in a fixed pivot, which moveable contact carrier can be operated by means of a breaker mechanism and is latched separately from the breaker mechanism by a locking lever. Thus, during rapid opening by means of a striker on the moveable contact carrier, a tension lever is moveable with respect to the contact compression force and thus the latch on the contact carrier is released. As such, the free ends of the locking lever and of the tension lever, which are mounted in the fixed pivot and which guide a compression spring, are joined together in an articulated manner, and support a sliding roller in the region of the articulated joint. The moveable contact carrier is provided with an essentially kidney-shaped slide plate in the direction of the switching operation, which slide plate forms a slip-in guide for the spring-tensioned sliding roller. A projection provided on the sliding plate in the proximity of the lock forms, together with the sliding part, a latching point. The angle of rotation of the tension lever is approximately the same as that of an angle of opening of the moveable contact carrier limited by a stop.\nApart from the fact that, in German Patent Application No. 31 33 285 C2, only a single finger contact system is described, this contact piece has disadvantages in that separate means are required to limit the opening movement caused by the short circuit as well as for collecting the kinetic energy occurring therewith, and that, with increasing wear of the contacts, the contact compression forces diminish."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical fiber adapter, and more particularly, to an optical fiber adapter with a shutter member.\n2. Description of the Related Art\nRecently, the optical fiber has been widely used as signal transmission media because of its high bandwidth and low loss nature. In order to transmit over a longer distance without the need of repeaters, it is common to use a high power diode laser to launch a laser beam into the optical fiber. However, the high power laser beam used to carry information is usually invisible. In other words, the human eyes will be unable to sense the high power laser beam if it leaves from the open end of a fiber cable. Thus, it is required to obstruct the high power laser beam to avoid hurting the eyes when the laser beam leaves from an optical fiber.\nReferring to FIG. 1, a conventional fiber adapter 100 includes a housing 110 having a plurality of side walls 160. The side walls 160 define a receiving recess 120. One of the side walls 160 is provided with a slot 130 for mating with the key 192 of the connector 190 when the connector 190 is inserted into the receiving recess 120. In addition, the outer surfaces of two opposite side walls 160 are provided with recesses 140 that a clip (not shown in the figure) can be disposed thereon to facilitate the adapter 100 to be mounted on a panel.\nIn general, the adapter 100 further has another set of side walls 160 that defines another receiving recess 120. The two receiving recesses 120 are opposite and can respectively receive a connector 190. Of course, the two receiving recesses 120 of the adapter 100 can be designed to mate with two different types of connectors. The fiber connector 190 is always attached to one end of a fiber cable 194 and a light beam can propagate down the fiber cable 194 and emit from the ferrule 196 of the connector 190. Likewise, a light beam can be coupled into the fiber cable 194 from the end face of the ferrule 196.\nWhen the adapter 100 is used to couple two connectors 190 together, the two connectors 190 are respectively inserted into the receiving recesses 120. The ferrules 196 of the connectors 190 thus slide into a hollow sleeve (not shown in the figure) and are brought into axial alignment and contact with each other. A light beam will be able to propagate from the fiber cable 194 of a connector 190 through the interface between the two ferrules 196 and then reach the fiber cable 194 of the other connector 190, and vice versa.\nWhen one connector 190 is disconnected from the adapter 100, the light beam originally propagating from the connector 190 still kept in the adapter 100 to the presently disconnected connector 190 will now leave the ferrule 196 and emit from the receiving recess 120. If the emitted light beam is high power and not obstructed, a lasting exposure to such light beam is harmful to people, particularly to the eyes. Thus, to avoid exposing to the high-power light beam, it is common to use a cap 180 to block up the unused receiving recess 120. This can obstruct the light beam and also prevent the receiving recess 120 from dust. If a connector 190 would like to be mated with the adapter 100, the cap 180 is required to be removed. However, the cap 180 is apt to get lost and it is still possible to expose the eyes to the light beam during mating.\nReferring to FIG. 2, a conventional fiber adapter 200 is generally the same as the adapter 100 but further includes a cover 250 pivotally connected to the housing 110. The cover 250 covers the receiving recess 120 in its closed position. A spring 260 can force the cover 250 to be pivoted to its closed position when the receiving recess 120 is not mated with a connector 190. Thus, the cover 250 is capable of obstructing the light beam emitted from the receiving recess 120 in its closed position. If a user would like to mate a connector 190 with the adapter 200, he is first required to lift the cover 250 from its closed position and then to insert the connector 190 into the receiving recess 120. Upon pulling out the connector 190, the cover 250 is pivoted to cover the receiving recess 120 through the spring 260. As a result, the user has no chance to expose to the high-power light beam. However, the construction of the adapter 200 is much more complex than that of the adapter 100.\nReferring to FIG. 3, a conventional protection cap 300 for the fiber adapter 100 includes a hollow housing 310 with two opposing openings and a cover 350 pivotally connected to the housing 310. The cover 350 is pivoted to cover one of the two openings through a spring 360. The protection cap 300 can be put on the adapter 100 and cover the outer surfaces of the side walls 160 and the receiving recess 120 of the adapter 100. When a user would like to mate a connector 190 with the adapter 200, he is required to lift the cover 350 from its closed position and then to insert the connector 190 into the receiving recess 120. Likewise, upon pulling out the connector 190, the cover 350 is pivoted to cover the receiving recess 120 by the spring 360. As a result, the user has no chance to expose to the high-power light beam emitted from the receiving recess 120. However, the structure of the protection cap 300 is somewhat complex and therefore it is not inexpensive. The introduction of the protection cap 300 to obstruct the light beam will cost much.\nAccordingly, there exists a need to provide a solution to solve the aforesaid problems."}
{"text": "The fabrication of integrated circuits on semiconductor wafers typically involves the creation of multiple, successive layers of materials, such as insulators, conductors, semiconductors, and so forth, on the semiconductor wafers. Each of the layers is normally formed by applying a photoresist layer over previously formed layers and then the photoresist layer is patterned. Portions of the photoresist layer not exposed during the patterning can be washed off and the layer can be formed using one of many desired techniques. When completed properly, the multiple layers combine to form functional integrated circuits. The alignment of the individual layers is crucial to creating properly formed structures in the semiconductor wafer. Misalignment of the layers can reduce the performance of the integrated circuitry if the misalignment is small (due to improper device geometries) and inoperable integrated circuitry if the misalignment is large (due to formation of improper electrical connections).\nMisalignment of the layers can arise when a photomask is used to pattern a layer and the photomask is not lined up properly with previously created layers on the semiconductor wafer. Misalignment can be due to mechanical shift errors, optical lens magnification errors, optical lens aberration errors, and so forth. Mechanical shift errors can be the result of shifts in the semiconductor wafer and/or the photomask during processing and optical lens magnification errors can be a result of magnification mismatches between different layers of the semiconductor wafer. Optical lens aberration errors can be the result of non-ideal characteristics of an optical lens being used, wherein light passing through the lens behaves differently depending upon the portion of the lens the light is passing through.\nOverlay targets have been used to allow for the alignment of the individual layers of a semiconductor wafer. After a photoresist layer has been patterned (and before the actual layer has been created) on top of the semiconductor wafer, which can have one or more targets, an optical system, such as a part of an overlay metrology tool, can capture images of the targets along with corresponding bullets (a part of the photoresist layer corresponding to the overlay target) in the photoresist layer and optical analysis algorithms can determine if the photoresist layer is aligned within specifications with respect to the semiconductor wafer. If the photoresist layer being created is determined to be misaligned, the photoresist layer can be removed, usually using a chemical wash, and a new photoresist layer can be applied and patterned and the optical processing can be repeated until the photoresist layer becomes aligned.\nOne disadvantage of the prior art is that the overlay targets can be used to detect mechanical shift errors and lens magnification errors. However, lens aberrations can also result in significant misalignment errors and the prior art overlay targets do not adequately capture lens aberration errors.\nA second disadvantage of the prior art is that the overlay targets do not exploit the advantages of using polarized light.\nYet another disadvantage of the prior art is that the overlay targets typically comprise the photoresist layer and a layer immediately beneath it. The use of adjacent layers can permit a sequential build-up of misalignment that, while between the adjacent layers is within specifications, can lead to an overall misalignment that exceeds specifications."}
{"text": "The invention relates to a transmission apparatus in connection with a pulling equipment intended for moving especially a cable or a corresponding product, the transmission apparatus comprising a shaft that is operated by a motor and that is connected to an output shaft of the transmission apparatus by means of cog wheels, said output shaft being connected to means that are in contact with the cable or a corresponding product and that move the cable, the motor being arranged to rotate the output shaft via the driven shaft and the cog wheels in the direction of motion of the cable or the like and, if necessary, to decelerate the rotation of the output shaft in the direction of motion of the cable or the like, the cog wheels being arranged to form a circle extending from the motor-driven shaft to the output shaft and correspondingly from the output shaft to the motor-driven shaft.\nIn addition to cables, the phrase xe2x80x9ca cable or a corresponding productxe2x80x9d refers in the present application to ropes, conductors and other products formed from strands, and to semi-finished articles of the products at different stages of production.\nWhen sheathed cables are produced, the sheath must be vulcanized and cooled after the forming stage. For this purpose, the processing line is provided with a vulcanizing tube that is followed by a cooling pipe in the direction of travel of the cable. A metal conductor to be coated is passed through the vulcanizing tube and the cooling pipe by means of a first pulling equipment placed before the extruder and a second pulling equipment located after the cooling pipe. During the vulcanization, the cable is supported in the horizontal extruder process entirely by the aforementioned pulling equipments or by the first pulling equipment, the cooling pipe and the second pulling equipment. The cable travels along a catenary trajectory at least in the vulcanizing tube, which means that at least the vulcanizing tube must be catenary in shape.\nDuring the vulcanization process, the cable must not get into contact with the hot vulcanizing tube, and therefore it is necessary that the cable is kept straight so that it does not touch the inner surface of the vulcanizing tube. If the cable touches the surface of the vulcanizing tube, the cable sheath will be damaged and the cable cannot be used for the intended purpose. The aforementioned matters set great requirements for the operation of the pulling equipments. A pulling equipment, for example a drive wheel, must grip the metal conductor without any slipping between the conductor and the drive wheel, the frame of the pulling equipment must be very sturdy and firmly attached to the bearing surface, the drive wheel must be firmly mounted in bearings to the frame in order to minimize the risk of vibration and oscillation, and for the same reason the transmission of the pulling equipment must operate as smoothly as possible without clearances. The clearance-free operation of the transmission of the pulling equipment is highly significant especially when loading takes place mainly in the direction of motion, as for example in the cable vulcanization process, where the motor may get into decelerating action in connection with certain cable sizes due to the tensile force of the cable. In other words, the motor may change its pole, whereupon the clearance of the gear transmission produces a tug, which easily results in harmful oscillation of the cable.\nIn prior art pulling equipments, rotation of the drive wheel is based on a transmission apparatus that is formed of a gear rim and a cog wheel and that is placed between an actuator and the drive wheel. In such arrangements, the drive wheel is provided with a gear rim that is coaxial with the axis of rotation of the drive wheel and that comprises either inside or outside gearing, and the pulling equipment is provided with a motor-driven cog wheel. The speed of the motor is controlled by means of a sensor placed in the vulcanizing tube and used to measure the location of the sheathed cable in the tube. The pulling equipment provided at the output end of the cooling pipe is usually run at a constant speed, and the location of the cable in the vulcanizing tube is adjusted by means of control of the rotational speed of the pulling equipment located before the extruder.\nIn practice, it has been noted, however, that despite very accurate machining of the gear teeth, vibrations occur in the pulling equipment in connection with the gear transmission described above, and the vibrations result in harmful oscillation of the cable in the vulcanizing tube. This is obviously a result of too big clearances between the teeth, the clearances becoming more evident with heavy loads. It should be noted that in the present cable vulcanization lines the distance between the pulling equipments may be as much as 200 meters, which means that a hanging cable causes considerable loading in the drive wheel and the teeth. A possible elliptical shape of the gear rim also causes variation in the clearances of the teeth and therefore in the steadiness of the pulling.\nIn another prior art arrangement, the drive wheel is replaced with several closely positioned rollers that are mounted in bearings in a substantially semicircular formation, and an inner endless rubber belt travels over the rollers. Outside the curved roller conveyor there is an outer endless rubber belt that passes over a sheave and presses against the roller conveyor in a semicircular form. The metal conductor to be coated travels between the inner and the outer belt to extruders. The speed of the metal conductor in the extruder is adjusted by decelerating the speed of the sheaves around which the inner belt travels. With this arrangement it is possible to avoid the drawbacks resulting from the clearances of the gear transmission, but on the other hand, the decelerating effect provided by means of the drive wheel is simultaneously lost.\nAnother essential drawback is that the inner rubber belt is subjected to a great pressure that is produced by the rollers and the outer rubber belt and that tends to lead to a breakdown of the inner rubber belt. Another drawback is that the outer rubber belt tends to slide away from the drive groove and therefore breaks.\nAnother example of prior art arrangements is a caterpillar-type pulling mechanism, which is arranged in connection with a drive wheel that is firmly mounted in bearings in the frame structure, the pulling mechanism being used to forward the metal conductor. On the one hand, such an arrangement improves for example the decelerating effect, but on the other hand, this kind of separate caterpillar equipment comprises clearances, which cause problems with vibration as described above.\nAn example of prior art arrangements is a pulling equipment where the force is conducted to the central shaft of the drive wheel by means of a worm gear. Such a gear is loaded with a high torque and therefore the equipment must be provided with a rather heavy construction, which also means that it has often rather great clearances that also cause the aforementioned problems with vibration.\nThe purpose of the present invention is to provide a transmission apparatus with which the prior art drawbacks, such as the vibration problems caused by the clearances of the transmission, can be eliminated. This is achieved with a transmission apparatus according to the invention, which is characterized in that the cog wheels are helical gearwheels and that one cog wheel shaft is arranged to be moved in its longitudinal direction when the rotation of the output shaft is being decelerated, whereupon during the movement of the cog wheel shaft, the cog wheels situated on the shaft are arranged to rotate, due to opposite helicity of the teeth of the cog wheels, the adjacent cog wheels that are in mesh with the first cog wheels, thus eliminating the free play occurring between the cog wheels during the deceleration.\nThe primary advantage of the invention is that the invention provides even transmission that does not cause harmful vibrations in the cable or a corresponding product. Another advantage of the invention is its simplicity, wherefore the start-up and use of the invention will be advantageous."}
{"text": "To respond a requirement of increase in withstanding voltage and decrease in on-resistance (or decrease in on-voltage) of a semiconductor device, a semiconductor device having the SJ structure has been developed, and in particular a semiconductor device having the SJ structure in both of the cell region and the terminal region is actively developed. The SJ structure is provided in not only in the cell region but also in the terminal region, thereby a depleted region (meaning an expansion area of a depletion layer when the semiconductor device is turned off) can be formed in a wide area of the terminal region, and consequently withstanding voltage of both of the cell region and the terminal region can be improved.\nThis type of semiconductor device is often formed by using a semiconductor stack in which a semiconductor lower-layer and a semiconductor intermediate-layer are stacked. In the semiconductor intermediate-layer, the SJ structure is formed. In many cases, when the semiconductor stack is seen in a plane view, the cell region in which a vertical semiconductor-switching-cell group is built is provided in a central side of the semiconductor stack, and the terminal region is provided around the periphery of the cell region.\nA plurality of vertical switching cells is formed in the cell region. For example, in the case that the vertical semiconductor switching cell is MOSFET (Metal Oxide Semiconductor Field Effect Transistor), the semiconductor lower-layer is called a drain layer, and a drain electrode is connected to the drain layer. When the vertical semiconductor switching cell is IGBT (Insulated Gate Bipolar Transistor), a collector layer in an opposite conduction type is provided on a back of the semiconductor lower-layer, and a collector electrode is connected to the collector layer. When a gate structure of the vertical semiconductor switching cell is in a planar type, a body region is formed in an upper region of the SJ structure, and a planar gate electrode is formed opposite to the body region. Alternatively, when the gate structure of the vertical semiconductor switching cell is in a trench type, a semiconductor upper-layer in a conduction type opposite to the semiconductor lower-layer is formed on a surface of the SJ structure, and a trench gate electrode that penetrates the semiconductor upper-layer is formed.\nAn avalanche resistance check of a semiconductor device is often carried out by an L-load surge resistance test and the like. In the L-load surge resistance test, excessive energy is supplied to the semiconductor device to forcibly induce breakdown. The breakdown is induced in a region of more than critical electric-field strength. Considering an area ratio between the cell region and the terminal region, breakdown is induced in a side of the cell region having a larger area, thereby avalanche energy per unit area can be reduced compared with a case that breakdown is induced in a side of the terminal region having a smaller area. Therefore, if breakdown is induced in the cell region side, local consumption of excessive avalanche energy can be suppressed, consequently occurrence of breakdown of the semiconductor device can be suppressed. To realize the phenomenon, withstanding voltage of the terminal region must be made high compared with standing voltage of the cell region, so that breakdown is induced dominantly in the cell region.\nJP-A-2003-273355, which corresponds to US 2005-0098826-A1 and U.S. Pat. No. 6,844,592, proposes a semiconductor device in which an insulating layer and a field plate are provided on a surface of the SJ structure in the terminal region. Furthermore, it proposes a structure in which thickness of the insulating layer is increased stepwise from the side of the cell region to the side of a region opposite to the cell region.\nWhen a combined structure of an insulating layer and a field plate is provided in the terminal region, an electric field in an upper region of the terminal region can be reduced. While the above prior art does not describe the following operation and effects in detail, when an insulating layer whose thickness is increased stepwise is used, an electric field near a boundary between the cell region and the terminal region, in which the electric field tends to be concentrated, can be reduced. It has been found from study of the inventors that an insulating layer which was adjusted to be thin is formed in the cell region side, thereby a significant effect of reducing the electric field near the boundary between the cell region and the terminal region is obtained. Thus, breakdown of the semiconductor device due to local concentration of the electric field can be avoided.\nHowever, as described before, in this type of semiconductor device, it is important that a relationship of “terminal region>cell region” is established between the withstanding voltage of the cell region and the withstanding voltage of the terminal region. The withstanding voltage of the cell region and that of the terminal region are mainly determined by height in a thickness direction of depleted regions formed in respective regions if the breakdown due to the local concentration of the electric field does not occur. In the semiconductor device of the above prior art, the height in the thickness direction of the depleted regions formed in both the cell region and the terminal region corresponds to height in the thickness direction of SJ structures formed in both. The height in the thickness direction of the cell region is equal to that of the terminal region. Therefore, the semiconductor device of the patent literature 1 has a limitation in that the withstanding voltage of the terminal region can be made equal to the withstanding voltage of the cell region at the maximum, or can not be made high compared with the withstanding voltage of the cell region."}
{"text": "Few of literatures and patents about the catalyst which can be used to produce ethylene by direct conversion of methane is disclosed; moreover, it is different in technical field from present invention. By conventional processes, large amounts of by-product such as carbon dioxide are produced, separation and removal thereof is difficult and environmental pollution is likely to be caused. Also, synthesis of hydrocarbon(s) by conventional dehydrogenation is conducted at relatively high temperature of about 1500.degree. to 1550.degree. C. through thermal or electric cracking reaction to cause some problems such as necessity of supply of high energy, expense of high temperature equipment (plan) as well as enormous loss of thermal energy expense, particulary severe corrosion of reactor.\nAs prior art describing synthesis of hydrocarbon by oxidative coupling or dehydrogenation reaction, there are U.S. Pat. Nos. 5,066,629, 5,068,486 and 5,118,654, Canadian patent no. 2016675 and Japanese patent nos. 04352730, 04368342."}
{"text": "This invention relates to a semiconductor laser and a method of producing the same which has recently rapidly extended its application to light sources for various types of electronic equipment and optical apparatus and is now greatly in demand.\nOne of the important functions of a semiconductor laser which is required when it is used as a coherent light source for electronic equipment or optical apparatus is oscillation at a single spot, or single transversal mode oscillation. To realize this, it is necessary to suppress the spread of laser light and confine the light by a concentration current flowing through the laser element near the active region in which the laser light propagates. This type of semiconductor laser is normally called the stripe-type semiconductor laser.\nA relatively simple method for the stripe type is to only restrict current to be narrow. The laser according to this method needs a high threshold value for realization of the single transversal mode oscillation. The stripe type which can oscillate at the lowest threshold is a buried stripe type semiconductor laser (normally called the BH laser). The production of this type of laser, however, requires that the crystal growth process be performed twice whereas with other lasers this process need normally be performed only once, that is, it is necessary that after the multilayered films including an active layer are grown, the buried region be removed and then a buried layer be again grown. In addition, it is somewhat technically difficult to produce this laser.\nIt is an object of this invention to provide a method of producing a semiconductor laser in which the buried-stripe structure that is necessary to oscillate in a single transversal mode at a low threshold value can be produced by a single crystal growth process."}
{"text": "The present disclosure relates to technologies for browsing, indexing, sharing, summarizing, and searching videos.\nIn recent years, two major technology advancements have created an important trend in everyday life: digital video applications technologies have become increasingly popular; the Internet has enabled people to be increasingly connected. Mobile technologies have improved from making simple phone calls to smart phones and tablets that are powerful computer devices equipped with high definition cameras and connected in 4G wireless networks. Many people now have two cameras in the front and back of their mobile devices always ready to video record of everything happening in their lives. People are willing to share the videos with their friends or the whole world. Digital video signals can now be streamed to mobile devices and computers over the Internet, and to Digital TV at homes. People also want to review old videos to find something or a place or someone that they met. The amount of video recording per user has increased from once a month or a week to several times per day, especially more during holidays, vacations, and other special events. The task of tagging all the videos has become a real challenge to many users. In the near future, people want to use videos as they have been text and images.\nA challenge in digital video applications today is a lack of video tools for managing videos in the same flexible way as we manage web and text content. Today video searching is conducted based on tag (metadata) that users associated to videos. The searches are manual and are text based, which requires time and effort from the user.\nMoreover, user text tagging is always under the criteria of the user, which can change over time for the same user, and are usually different from the criteria of the user's family and friends with whom the user's videos are to be shared. Thus searching the even tagged videos can become a difficult task for the same user over time and for different users.\nIt is very difficult and time consuming to use current technologies to browse, organize, share, and search the long video files stored in smart phones, tablets, and computers.\nSome products have made initial attempt to address this problem by recognizing faces of people in the video and tag them as metadata. The tags can be used to search people in videos. But these products have limited applications because not all events involve people or human faces and even a person is recognized, more detailed classifications are often needed to recognize the nature of the events.\nSome approaches select a few frames that are called “key frames” which are a set of pictures or slides that capture key moments in a video. This “key frames” can significantly reduce file size from a video to a set of pictures, but it pays a heavy price because movement and audio, the essential properties of a video, are lost in the process.\nThere is therefore still a need to provide simple and effective tools for people to browse, index, share, and search videos in a similar fashion as people manage text and web content."}
{"text": "In recent years, various kinds of electronic devices, such as a tablet terminal, a PDA and a smartphone, have been developed. Most of these electronic devices include touch-screen displays for facilitating input operations by users.\nBy touching a menu or an object, which is displayed on the touch-screen display, by a finger or the like, the user can instruct an electronic device to execute a function which is associated with the menu or object.\nHowever, most of existing electronic devices with touch-screen displays are consumer products which are designed to enhance operability on various media data such as video and music, and are not necessarily suitable for use in a business situation such as a meeting, a business negotiation or product development. Thus, in business situations, paper-based pocket notebooks have still been widely used.\nRecently, a technique for determining a relationship between a tablet terminal and the seating position of a user, based on the direction of writing of handwritten characters, has also been developed.\nIn the meantime, in general, most of tablet terminals have a function of automatically rotating the direction of a screen image in accordance with the direction of the tablet terminal.\nHowever, in a handwriting application which handles a handwritten document, if the direction of the screen image of the application is automatically rotated in accordance with the direction of the tablet terminal, there may be a case in which a feeling of use, such as a feeling of using a real paper-based pocket notebook, cannot be obtained. The reason for this is that if the direction of the screen image is automatically rotated, it would become difficult to handwrite characters, etc. in a free direction on the display screen of the handwriting application.\nOn the other hand, the handwriting application is required to also have a function as a digital tool for handling a handwritten document as digital data."}
{"text": "Conventionally, a refrigerating apparatus configured to cool an inside of a compartment such as a container used for, e.g., marine transportation or a trailer used for, e.g., ground transportation has been known. For example, Patent Document 1 discloses the refrigerating apparatus (trailer refrigerating apparatus) of this type. The refrigerating apparatus includes a casing body attached to an opening of a trailer of a refrigerated vehicle, and a refrigeration cycle unit supported by the casing body. The refrigeration cycle unit forms a refrigerant circuit configured to perform a vapor compression refrigeration cycle. Specifically, the refrigeration cycle unit includes a compressor, a radiator, an expansion valve, and an evaporator. The compressor, the radiator, and the expansion valve are supported by the casing body so as to be positioned outside a cooling compartment (i.e., outside a compartment). The evaporator is supported by the casing body so as to be positioned inside the cooling compartment (i.e., inside the compartment).\nAn air path is formed inside the cooling compartment (trailer) such that air circulates through the evaporator. On the other hand, when refrigerant circulating through the refrigerant circuit passes through the evaporator, the refrigerant is evaporated by absorbing heat from in-compartment air. In such a manner, the in-compartment air of the trailer is cooled to store food etc."}
{"text": "Embodiments of the present invention relate generally to the processing of integrated circuit elements and in particular to the processing of contact pads of integrated circuit elements.\nContact pads (also referred to as bond pads or short pads) are commonly used to make external electrical contact with an integrated circuit (IC) element or device, e.g., with a chip."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device. In particular, the present invention relates to a semiconductor device functioning as a driver circuit for driving a power device for power supply.\n2. Description of the Related Art\nA power device is an element used to drive a load such as a motor. The power device can control supply of a large current to a load to perform the supply intermittently by a switching operation.\nAs an example of the power device, a power transistor such as an insulated gate bipolar transistor (IGBT) can be given. The power transistor has a large gate capacitance compared to other transistors.\nDriving of the power transistor is controlled by pulse width modulation (PWM). PWM control is performed in accordance with a PWM signal output from a microcomputer or the like. The voltage of the PWM signal is too low to directly drive the power transistor having a large gate capacitance. Therefore, the PWM signal needs to be converted into a high-voltage signal and supplied to the power transistor.\nA driver circuit for converting the PWM signal to a high-voltage signal includes a transistor including silicon. For example, Patent Document 1 discloses a structure of a semiconductor device in which an n-channel transistor and a p-channel transistor are provided over a silicon substrate and the on/off state of IGBT is controlled."}
{"text": "1. Field\nThe present disclosure relates to a liquid crystal display. More particularly, the present disclosure relates to a liquid crystal display having a backlight assembly.\n2. Description of the Related Art\nIn recent years, various display devices have been developed to be applicable to a portable information processing device, such as a tablet computer, a notebook computer, etc. In addition to ways to improve a display quality of the display devices, various approaches have been studied to reduce the display devices in size and weight."}
{"text": "1. Field of the Invention\nThe present invention refers to a coating device for the deposition of thin films on especially large-area substrates, like architectural glass or display panels. Furthermore, the present invention refers to a method for coating such substrates.\n2. Prior Art\nCoating devices for coating of especially thin layers on substrates like glass are widely used in industrial applications, e.g. depositing of heat insulation layers on architectural glass or conductive or semiconducting layers on display panels or photovoltaic units.\nIn order to apply such thin films on large-area substrates, devices and methods are known from prior art. However, a permanent need for improvement exists with respect to an increase of homogeneity of layer thickness or other properties of the deposited layers over the whole coated substrate as well as an enhancement of efficiency of coating.\nSince it is always difficult to find an appropriate comprise to achieve these diverging needs, an improvement in this respect is especially valuable.\nFrom prior art sputter processes are known for carrying out coating processes mentioned above.\nAs can be seen from FIG. 9 showing a prior art apparatus, sputter devices have an electrode arrangement 401, 402 within a process chamber 201, 202 in order to ignite a plasma and accelerate ions of the plasma towards a target in front of a cathode of the electrode arrangement 401, 402. Magnetron cathodes comprise a magnet arrangement, which is disposed behind the target so as to lead to a confinement of electrons in front of the target and thus to a plasma concentration in this area. Thus, a high sputter yield is achieved due to the high concentration of ions impinging on the target. The material of the target is atomized and deposits on the substrate 5, 6 being disposed opposing to the magnetron cathode.\nIn order to achieve a high through-put of the device, a double-chamber design can be used with two process chambers 201, 202 arranged side by side to each other. Such a double-chamber 200 enables a production of the process chambers with reduced effort and allows a space saving design of the sputter apparatus. According to the embodiment shown in FIG. 9, a buffer chamber 100 for storing the substrates to be coated and a lock section 300 having two locks 301 and 301 for feeding the substrates into both process chambers 201, 202 are additionally provided for.\nDE 41 26 236 C2 also describes a device for a sputter process. This device comprises a rotating magnetron cathode having a hollow cylinder target, which rotates in order to use the whole target surface for sputtering.\nAccording to DE 41 26 236 C2 two magnet arrangements are disposed in the hollow cylinder target so as to form two plasma areas on opposing sides of the cylindrical magnetron electrode. One plasma is used for coating the substrate, the second plasma generated by the second magnet arrangement is used for cleaning of the rotating target. However, it is also possible to simultaneously deposit layers on two substrates being located at the opposing sides of the cathode. This device can however only be used for small-area coatings or a dynamic coating process wherein the substrates are passing the single sputter source.\nA similar device is disclosed in U.S. Pat. No. 5,158,660. However, instead of a magnet arrangement inside the cylindrical electrode, a magnet arrangement disposed outside the target or cathode is used. However, this device of U.S. Pat. No. 5,158,660 is also a single coating source with small lateral dimensions for small-area coatings or dynamic coating processes.\nIn order to achieve large-area coatings with homogenous layer thicknesses, EP 15 94 153 A1 suggests to provide a plurality of cathodes in order to allow stationary coating of large-area substrate facing the plurality of cathodes. However, in order to be sure that at the margins of the substrates the same layer thickness is achieved as in the central area, the cathodes of EP 15 93 153 A1 are arranged at a curved surface. Accordingly, it is not possible to use cathodes described in DE 41 26 236 C2 or U.S. Pat. No. 5,158,660 having deposition areas on two sides of the cathodes, since the curved arrangement of the cathode does not allow homogenous twin-deposition from two sides of the cathode.\nOther devices, e.g. described in WO 2007/051105 A2 or JP 2004 027 272 A, deal with the above-mentioned problems by providing movable magnet arrangements for magnetron electrodes. Although, this leads to an improved material usage of target material and allows for adapted deposition areas leading to an enhancement of layer homogeneity, improvements with respect to layer homogeneity of large-area depositions and improvements with respect to efficiency are still necessary."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of semiconductor device manufacturing and, more particularly, to a method and apparatus for determining grid dimensions using scatterometry.\n2. Description of the Related Art\nThere is a constant drive within the semiconductor industry to increase the quality, reliability and throughput of integrated circuit devices, e.g., microprocessors, memory devices, and the like. This drive is fueled by consumer demands for higher quality computers and electronic devices that operate more reliably. These demands have resulted in a continual improvement in the manufacture of semiconductor devices, e.g., transistors, as well as in the manufacture of integrated circuit devices incorporating such transistors. Additionally, reducing the defects in the manufacture of the components of a typical transistor also lowers the overall cost per transistor as well as the cost of integrated circuit devices incorporating such transistors.\nThe technologies underlying semiconductor processing tools have attracted increased attention over the last several years, resulting in substantial refinements. However, despite the advances made in this area, many of the processing tools that are currently commercially available suffer certain deficiencies. In particular, such tools often lack advanced process data monitoring capabilities, such as the ability to provide historical parametric data in a user-friendly format, as well as event logging, real-time graphical display of both current processing parameters and the processing parameters of the entire run, and remote, i.e., local site and worldwide, monitoring. These deficiencies can engender nonoptimal control of critical processing parameters, such as throughput, accuracy, stability and repeatability, processing temperatures, mechanical tool parameters, and the like. This variability manifests itself as within-run disparities, run-to-run disparities and tool-to-tool disparities that can propagate into deviations in product quality and performance, whereas an ideal monitoring and diagnostics system for such tools would provide a means of monitoring this variability, as well as providing means for optimizing control of critical parameters.\nSemiconductor devices are manufactured from wafers of a substrate material. Layers of materials are added, removed, and/or treated during fabrication to create the electrical circuits that make up the device. The fabrication essentially comprises four basic operations. Although there are only four basic operations, they can be combined in hundreds of different ways, depending upon the particular fabrication process.\nThe four operations typically used in the manufacture of semiconductor devices are:                layering, or adding thin layers of various materials to a wafer from which a semiconductor device is produced;        patterning, or removing selected portions of added layers;        doping, or placing specific amounts of dopants in the wafer surface through openings in the added layers; and        heat treatment, or heating and cooling the materials to produce desired effects in the processed wafer.        \nAmong the important aspects in semiconductor device manufacturing are rapid thermal annealing (RTA) control, chemical-mechanical polishing (CMP) control, etch control, and photolithography control. As technology advances facilitate smaller critical dimensions for semiconductor devices, the need for reduction of errors increases dramatically. Proper formation of sub-sections within a semiconductor device is an important factor in ensuring proper performance of the manufactured semiconductor device. Critical dimensions of the sub-sections generally have to be within a predetermined acceptable margin of error for semiconductor devices to be within acceptable manufacturing quality.\nGenerally, most features on a semiconductor device are formed by depositing layers of material (e.g., conductive or insulative) and patterning the process layers using photolithography and etch processes. The various process layers used for forming the features have many specialized functions. Certain layers are used to form conductive features, others are used to form insulating features, and still others are intermediate layers used to enhance the functionality of the processing steps used to pattern and form the functional layers.\nThere are many variables that affect the accuracy and repeatability of the photolithography and etch processes used to form features from the process layers. Typical metrology data collection for measuring the efficacy of the photolithography and etch processes do not provide data that is sufficiently accurate and timely to facilitate run-to-run control of such processes. Certain techniques, such as scanning electron microscope (SEM) analysis, may be used to generate two-dimensional data, such as a critical width dimension, but they cannot be used to characterize the entire feature. For example, to generate a sidewall angle measurement, an important measure of quality, a destructive cross-section SEM analysis is required. Destructive tests are expensive, as they require destruction of the wafer, and thus, it is impractical to perform such tests at a frequency that would allow run-to-run control.\nThe present invention is directed to overcoming, or at least reducing the effects of, one or more of the problems set forth above."}
{"text": "In a production process of an aluminum alloy wheel, spraying is a very important procedure, and the procedure is carried out to achieve a better appearance effect for the wheel. However, the procedure contains numerous steps, after spraying bottom powder and colored paint on the wheel for the first time, and before carrying out the next step, bottom powder on a ring flange of the wheel and at a center hole must be removed; at this time, if the bottom powder is removed manually, the efficiency is absolutely very low, and the effect is not good; and therefore, an automatic device is needed to quickly remove the bottom powder online."}
{"text": "1. Field of the Invention\nThe present invention relates to sensors capable of detecting a gas in a fluid and, more particularly, to sensors capable of measuring the presence and the amount of oxygen (O.sub.2).\n2. Description of the Prior Art\nAnalytical methods for the determination of O.sub.2 in a fluid include the amperometric and galvanic methods. These methods are quite rapid, simple in operation, and are specially suited for determining O.sub.2 in either its gaseous state or as it may be found dissolved in a liquid.\nIn the amperometric method, a sensor such as that shown in FIG. 1 is employed. In such a sensor, generally indicated as 10, body 12 is provided having a reservoir 14 therein containing a center post 16 for supporting a gold cathode 18. A membrane holder 20, cap 22, and O-ring 23 are provided to hold a polymeric membrane material 24, such as polytetrafluoroethylene of a thickness of from 0.00025 inches to 0.002 inches, stretched over the gold cathode thereby completely enclosing reservoir 14. Reservoir 14 is then filled with an electrolyte 26, typically a 5% KCl solution, either buffered or unbuffered. When the sensor is to be used in an application where temperatures in excess of 100.degree. F. are generally encountered, the electrolyte is generally mixed with typical antifreeze compounds such as ethylene glycol. Additionally, an anode 28, typically silver, is disposed within reservoir 14 in contact with electrolyte 26. A potential of about 750 millivolts is applied between the anode 28 and the cathode 18 by means 29 connected thereto. In stretching the membrane 24 over cathode 18, a very minimal amount of electrolyte 26 is contained between gold cathode 18 and membrane 24. As sample fluid is brought in contact with membrane 24, O.sub.2 diffuses through membrane 24 to contact the gold cathode 18 in the presence of the electrolyte 26. A current flow results which is linear with the partial pressure of O.sub.2 being sampled. Thus, this current can be measured and correlated to the amount of O.sub.2 in the sample.\nIn the galvanic method, a sensor such as that shown in FIG. 2 is employed, wherein corresponding numbers indicate corresponding parts. The sensor employed in the galvanic method typically has a silver cathode 18', a lead anode 28' and a readout device 29', e.g., an ammeter.\nWhile these methods are satisfactory in many applications, they suffer drastically in atmospheres of high carbon dioxide (CO.sub.2) such as encountered in monitoring automobile exhaust or stack gas emissions. Such sensors as that of FIG. 1 and FIG. 2, as described above, become less sensitive to O.sub.2 upon even brief exposure to high concentrations of CO.sub.2 and may take several hours to recover so as to indicate the proper value of O.sub.2. The response of such sensors is highly dependent on the pH of the electrolyte at the interface between the gas diffusing through the membrane and the cathode. At low, or acid, pH levels the response is low. At high, or base, pH levels the response is higher. Typical responses by a prior art sensor are shown in FIG. 3 and FIG. 4. In FIG. 3, the sensor was first exposed to ambient air containing approximately 21% O.sub.2. It was then exposed to pure nitrogen (N.sub.2). FIG. 3 shows the response which was both expected and achieved. Upon exposure to N.sub.2, the response dropped to the zero line. Upon exposure to air, the response climbed to the level indicating approximately 21% O.sub.2. This cycle was repeatable without problem. Referring now to FIG. 4, the expected and actual response of a prior art sensor is shown when the sensor was exposed to ambient air and then exposed to a mixture of 15% CO.sub.2 plus 3% O.sub.2 and the balance nitrogen. When exposed to the mixture, the expected response is for the output to drop to the 3% level, being an indication of the 3% O.sub.2 content of the mixture. Upon exposure to ambient air, it is expected that the output will climb to the 21% O.sub.2 level of the ambient air. The actual response, however, was not as anticipated. When the sensor was exposed to the mixture, the response fell to the expected 3% O.sub.2 level. When the sensor was subsequently exposed once again to the ambient air sample, the output overshot the 21% level, then reversed and undershot the 21% level, and then slowly approached the 21% level asymptotically. It was found that the recovery period required for the output to attain the actual 21% level varied depending both on the duration of exposure to CO.sub.2 and the amount of CO.sub.2 in the sample. For example, when using such a prior art sensor on an automobile exhaust, an exposure for a period one minute to the exhaust gases resulted in a recovery period on the order of two to three hours before an accurate ambient response could be attained.\nThis phenomenon is a result of the small volume of electrolyte trapped adjacent to the gold cathode by the membrane. This is typically on the order of 1 microliter. As previously mentioned, the response of such a cell is dependent on the pH of the electrolyte. When CO.sub.2 is introduced, carbonic acid is formed which, when mixed with such an extremely small volume of basic electrolyte, results in a change of the pH of the electrolyte adjacent the gold cathode. In the typical electrolyte having a pH of approximately 13.5, the introduction of carbonic acid having a pH in the order of 4.5 results in a change of pH of the electrolyte trapped adjacent the gold cathode to a level of approximately 9. The response of the electrode will be correspondingly reduced until such time as the pH can attain its normal value by diffusion of normal electrolyte into the space between the membrane. The initial overshoot observed is, presumably, caused by the sudden change in pH and the unsettling of the electrolyte in the cathode area.\nAttempts at improving the performance of O.sub.2 sensors are not new in the art. It is well known that changes in pH of the electrolyte adjacent the cathode will change the response of the electrode. On the other hand, it is known that the thickness of the electrolyte in this same area affects the sensitivity of the electrode to oxygen. Thus, ideally, the spacing between the membrane and the cathode is kept minimal while means are provided for allowing the free movement of the electrolyte through the space. Thus, in the prior art, it has been suggested to roughen the surface of the gold cathode, provide channels therein for the movement of electrolyte, and depose porous materials between the membrane and the cathode to provide channels for the movement of the electrolyte.\nSuch prior art suggestions have resulted in O.sub.2 sensors of marginal sensitivity and poor response times for certain applications. In particular, automobile exhaust analysis and flue gas analysis provide environments imposing restrictions beyond the capabilities of prior art O.sub.2 sensors employing such techniques. In the field of automobile exhaust gas analysis, the ability to cycle and recover at rapid rates is imperative in \"assembly line\" type testing environments.\nAccordingly, it would be highly desirable to have an electrochemical O.sub.2 sensor capable of yielding accurate and reproducible date in applications wherein high CO.sub.2 levels are encountered at both ambient, or low, and elevated temperatures. Such a sensor would be highly useful for monitoring the O.sub.2 content of automobile exhaust and stack gas emissions."}
{"text": "1. Field of the Invention\nThe present invention relates to mounting assemblies and, more particularly, to a mounting assembly for storage devices.\n2. Description of Related Art\nA computer, such as a typical desktop computer, a tower computer, a server, and the like, usually includes storage devices, such as hard disk drives (HDDs). In using a computer, especially a server, different kinds of HDDs may be needed, each having a different size. Some standard HDD sizes include 2.5 inches, 3.5 inches, and so on. Thus, mounting assemblies capable of mounting different HDDs to a computer are needed."}
{"text": "Prior microwave electrical devices include resistors, series shunt elements, and solid-state devices such as Schottky diodes, PIN diodes, varactor diodes, and transistors. Another prior microwave device is a bridge quad (also referred to as a sampling bridge) that includes four diodes in one package. Another prior microwave electrical device that has four diodes is a ring quad, also called a quad-ring.\nOne type of prior microwave device has beam leads on the top of the device to provide connections to external circuitry. The beam leads are structurally supported by diodes that make up the device. The beam leads also provide interconnections between terminals of the diodes. For example, on certain prior ring quads, the beam leads provide interconnections between the anodes and cathodes of the respective diodes.\nFor another type of prior microwave electrical device, metal interconnections reside on one side of the device and metalized layers for external connections reside on the other side of the device. This helps to make the device less fragile. Prior microwave devices using this alternative configuration include ring quads, bridge quads, and diodes.\nBoth of the prior types of microwave devices discussed above typically use glass to hold components of the devices together. For example, certain prior ring quads typically use glass to hold four diode blocks together. Metallic interconnections for those prior microwave devices typically provide only minimal support. Thus, for certain types of prior microwave devices, glass is the material that provides major structural support.\nNevertheless, the structural support provided by the glass typically is only for a portion of a prior device assembly process. During the final stages of device fabrication, prior microwave devices are typically either (1) completely encased in glass to form a hermetic seal or (2) encapsulated in epoxy.\nThe way glass is used for certain prior microwave devices has several disadvantages. One disadvantage is that the glass makes the prior microwave devices relatively fragile. Because of this fragility, those prior microwave devices have relatively low maximum tear strengths. This means that certain prior microwave devices can be destroyed relatively easily.\nAnother disadvantage associated with certain prior microwave devices is that assembly operations are relatively tedious and delicate given that beam leads of certain prior microwave devices are not covered by the glass layer. Capacitance, inductance, and balance problems can arise from improper bonding of the beam leads. The exposed and relatively large beam leads during the fabrication process are also prone to bending, which can result in capacitance, inductance, and balance problems. Dynamic resistance problems can also arise. Special care is typically required for the beam leads to seat properly during packaging.\nAnother disadvantage associated with glass is that air bubbles are typically formed inside the glass layer during fabrication. When the liquified glass is deposited during certain prior fabrication processes, air often enters the liquified glass and forms air bubbles inside the glass. The air bubbles inside the glass typically weaken the glass in a nonuniform way from point to point and from device to device. This results in differences in strength between prior microwave devices from device to device.\nAnother disadvantage with certain prior microwave devices is that glass typically is liquified as part of the fabrication process. This requires relatively high temperatures.\nOne type of prior integrated circuit is typically attached to a package substrate when being packaged to a lead package. The leads of the prior integrated circuit are typically bonded to respective leads of the lead package.\nOne disadvantage of such prior integrated circuit is that the leads of the prior integrated circuit tend to have relatively high lead inductance that typically results in electrical problems with respect to the circuit. Another disadvantage is that the packaging of the prior integrated circuit is typically relatively expensive. Another disadvantage is that the leads need to be spaced apart far enough to permit the use of wire bonding machinery.\nAnother type of prior integrated circuit employs beam leads formed on the top of the prior integrated circuit for attachment to a next level of interconnection. In addition, the beam leads provide connections to external circuitry. These beam leads, however, are typically relatively fragile and difficult to handle.\nAnother type of prior integrated circuit uses bumps on the die or substrate or both, or uses solder balls, to attach the chip directly to a substrate in a \"flip chip\" fashion. One disadvantage with these approaches is that there may be different thermal expansions with respect to the chip and the substrate, which in some cases can lead to failure of the solder connection. Moreover, such arrangements may result in the chip being easily damaged during handling."}
{"text": "This invention relates broadly to the structure of macromolecules and more specifically to nuclear magnetic resonance (NMR) methods for preferentially observing exposed positions of a macromolecule.\nThe three-dimensional structure of macromolecules such as proteins can provide insight into their function in normal and diseased cells. For example, when certain compounds or “ligands” bind to a protein, the 3-D structure of the protein can help identify the ligand's binding sites on the protein. Atomic-resolution structures can allow researchers to understand the interaction between ligands and macromolecules and to design and select new ligands that may serve as therapeutic drugs.\nNuclear magnetic resonance (NMR) spectroscopy has determined the structure of many proteins at atomic resolution. Nevertheless, current NMR methods for protein structure determination are extremely time-consuming because they require painstaking deduction of the positions of thousands of individual atoms. As attention focuses on larger and more complex proteins and the number of atoms increases, the data generated by NMR methods can become overwhelming, making it virtually impossible to determine the structure of the protein.\nThus, there is a need for methods that can provide information about large macromolecular structures in a way that is useful for understanding ligand binding. The present invention satisfies this need and provides related advantages as well."}
{"text": "Display devices are used for presentation of video to a viewer. The video to be presented may be transmitted to the display device from a source unit in the form of a video signal. A display device receiving a video signal will render video output corresponding to data conveyed by the video signal. However, if the display device improperly interprets the video signal, which could occur, e.g., if the video signal includes excessive noise, the corresponding video output may be disrupted or compromised."}
{"text": "This invention is directed to bitumen compositions, which are prepared from bitumen, polymers such as copolymers of styrene and a conjugated-diene, and defined amounts of crosslinking agents such as sulfur. The bitumen compositions described herein are useful in industrial applications; such as in hot mix asphalts useful in preparing aggregates for road paving.\nThe use of bitumen (asphalt) compositions in preparing aggregate compositions (bitumen+rock) useful as road paving material is complicated by at least three factors, each of which imposes a serious impediment to providing an acceptable product. First, the bitumen compositions must meet certain performance criteria or specifications in order to be considered useful for road paving. For example, to ensure acceptable performance, state and federal agencies issue specifications for various bitumen applications including specifications for use as road pavement. Current Federal Highway Administration specifications designate a bitumen (asphalt) product, for example, AC-20R as meeting defined parameters relating to properties such as viscosity, toughness, tenacity and ductility (see Table 1). Each of these parameters define a critical feature of the bitumen composition, and compositions failing to meet one or more of these parameters will render that composition unacceptable for use as road pavement material.\nConventional bitumen compositions frequently cannot meet all of the requirements of a particular specification simultaneously and, if these specifications are not met, damage to the resulting road can occur, including permanent deformation, thermally induced cracking and flexural fatigue. This damage greatly reduces the effective life of paved roads.\nIn this regard, it has long been recognized that the properties of conventional bitumen compositions can be modified by the addition of other substances, such as polymers. A wide variety of polymers have been used as additives in bitumen compositions. For example, copolymers derived from styrene and conjugated dienes, such as butadiene or isoprene, are particularly useful, since these copolymers have good solubility in bitumen compositions and the resulting modified-bitumen compositions have good rheological properties.\nIt is also known that the stability of polymer-bitumen compositions can be increased by the addition of crosslinking agents such as sulfur, frequently in the form of elemental sulfur. It is believed that the sulfur chemically couples the polymer and the bitumen through sulfide and/or polysulfide bonds. The addition of extraneous sulfur is required to produce the improved stability, even though bitumens naturally contain varying amounts of native sulfur.\nThus, U.S. Pat. No. 4,145,322, issued Mar. 20, 1979 to Maldonado et al., discloses a process for preparing a bitumen-polymer composition consisting of mixing a bitumen, at 266°–446° F. (130°–230° C.), with 2 to 20% by weight of a block copolymer, having an average molecular weight between 30,000 and 300,000, with the theoretical formula Sx-By, in which S corresponds to styrene structure groups and B corresponds to conjugated diene structure groups, and x and y are integers. The resulting mixture is stirred for at least two hours, and then 0.1 to 3% by weight of sulfur relative to the bitumen is added and the mixture agitated for at least 20 minutes. The preferred quantity of added sulfur cited in this patent is 0.1 to 1.5% by weight with respect to the bitumen. The resulting bitumen-polymer composition is used for road coating, industrial coating, or other industrial applications.\nSimilarly, U.S. Pat. No. 4,130,516, issued Dec. 19, 1978 to Gagle et al., discloses an asphalt (bitumen) polymer composition obtained by hot-blending asphalt with 3 to 7% by weight of elemental sulfur and 0.5 to 1.5% by weight of a natural or synthetic rubber, preferably a linear, random butadiene/styrene copolymer. U.S. Pat. No. 3,803,066, issued Apr. 9, 1974 to Petrossi, also discloses a process for preparing a rubber-modified bitumen by blending rubber, either natural or synthetic, such as styrene/butadiene rubber, with bitumen at 293°–365° F. (145°–185° C.), in an amount up to 10% by weight based on the bitumen, then adjusting the temperature to 257°–320° F. (125°–160° C.), and intimately blending into the mix an amount of sulfur such that the weight ratio of sulfur to rubber is between 0.3 and 0.9. A catalytic quantity of a free-radical vulcanization-accelerator is then added to effect vulcanization. This patent recites the critical nature of the sulfur to rubber ratio, and teaches that weight ratios of sulfur to rubber of less than 0.3 gives modified bitumen of inferior quality.\nAlthough polymer-modified bitumen compositions are known, these previously described compositions are not necessarily useful for road paving applications. For example, mixing NorthWest paving asphalt having an initial viscosity of 682 poise at 140° F. (60° C.) with 3.6 weight percent Kraton®-4141, a commercially available styrene-butadiene tri-block copolymer which contains 29 weight percent plasticizer oil, and 0.25% sulfur gives a modified-asphalt composition with a viscosity of 15,000 poise at 140° F. (60° C.). This viscosity, however, greatly exceeds the acceptable viscosity range set by the widely used AC-20R specification for paving asphalt. This specification, issued by the Federal Highway Administration, requires bitumen compositions to have a viscosity in the range of 1600–2400 poise at 140° F. (60° C.). Thus, the modified bitumen compositions produced by the procedures of U.S. Pat. No. 4,145,322 using Kraton®-4141 would be unacceptable for use in road paving under the AC-20R specification.\nThe second factor complicating the use of bitumen compositions concerns the viscosity stability of such compositions under storage conditions. In this regard, bitumen compositions are frequently stored for up to 7 days or more before being used and, in some cases, the viscosity of the composition can increase so much that the bitumen composition is unusable for its intended purpose. On the other hand, a storage stable bitumen composition would provide for only minimal viscosity increases and, accordingly, after storage it can still be employed for its intended purpose.\nThe third factor complicating the use of bitumen compositions concerns the use of volatile solvents in such compositions. Specifically, while such solvents have been heretofore proposed as a means to fluidize bitumen-polymer compositions containing relatively small amounts of sulfur which compositions are designed as coatings (Maldonado et al., U.S. Pat. No. 4,242,246), environmental concerns restrict the use of volatile solvents in such compositions. Moreover, the use of large amounts of volatile solvents in bitumen compositions may lower the viscosity of the resulting composition so that it no longer meets viscosity specifications designated for road paving applications. In addition to the volatile components, reduction of other emissions during asphalt applications becomes a target. For example, it is desirable to reduce the amount of sulfur compounds that are emitted during asphalt applications.\nAsphaltic concrete, typically including asphalt and aggregate, asphalt compositions for resurfacing asphaltic concrete, and similar asphalt compositions must exhibit a certain number of specific mechanical properties to enable their use in various fields of application, especially when the asphalts are used as binders for superficial coats (road surfacing), as asphalt emulsions, or in industrial applications. (The term “asphalt” is used herein interchangeably with “bitumen.” Asphaltic concrete is asphalt used as a binder with appropriate aggregate added, typically for use in roadways.) The use of asphalt or asphalt emulsion binders either in maintenance facings as a surface coat or as a very thin bituminous mix, or as a thicker structural layer of bituminous mix in asphaltic concrete, is enhanced if these binders possess the requisite properties such as desirable levels of elasticity and plasticity.\nPreviously, various polymers have been added to asphalts to improve physical and mechanical performance properties. Polymer-modified asphalts are routinely used in the road construction/maintenance and roofing industries. Conventional asphalts often do not retain sufficient elasticity in use and, also, exhibit a plasticity range that is too narrow for use in many modern applications such as road construction. It is known that the characteristics of road asphalts and the like can be greatly improved by incorporating into them an elastomeric-type polymer which may be one such as butyl, polybutadiene, polyisoprene or polyisobutene rubber, ethylene/vinyl acetate copolymer, polyacrylate, polymethacrylate, polychloroprene, polynorbornene, ethylene/propylene/diene (EPDM) terpolymer and advantageously a random or block copolymer of styrene and a conjugated diene. The modified asphalts thus obtained commonly are referred to variously as bitumen/polymer binders or asphalt/polymer mixes. Modified asphalts and asphalt emulsions typically are produced utilizing styrene/butadiene based polymers, and typically have raised softening point, increased viscoelasticity, enhanced force under strain, enhanced strain recovery, and improved low temperature strain characteristics.\nThe bituminous binders, even of the bitumen/polymer type, which are employed at the present time in road applications often do not have the optimum characteristics at low enough polymer concentrations to consistently meet the increasing structural and workability requirements imposed on roadway structures and their construction. In order to achieve a given level of modified asphalt performance, various polymers are added at some prescribed concentration.\nCurrent practice is to add the desired level of a single polymer, sometimes along with a reactant that promotes cross-linking of the polymer molecules until the desired asphalt properties are met. This reactant typically is sulfur in a form suitable for reacting. Such current processes are discussed in various patents such as U.S. Pat. No. 4,145,322 (Maldonado); U.S. Pat. No. 5,371,121 (Bellamy); and U.S. Pat. No. 5,382,612 (Chauerat), all of which are hereby incorporated by reference.\nOther reactants and combinations of reactants are disclosed in U.S. Pat. No. 5,104,916 to Trinh et al., U.S. Pat. No. 3,992,340 to Bonitz and U.S. Pat. No. 4,129,542 to Matheson et. al. Common to these patents is the use of tetramethylthiuram disulfide (TMTDS) and/or stearic acid in significant concentration. Both of these compounds add raw material costs and handling difficulties the present invention avoids. For example, stearic acid, when airborne is a combustible dust. TMTDS is moderately toxic by ingestion but is highly toxic if inhaled. Further, TMTDS has been shown to break down into n-nitroso-dimethylamine in the human stomach, a compound that has been shown to be carcinogenic in test animals. Surprisingly, the present invention achieves the required asphalt properties without the use of TMTDS or stearic acid.\nThe cost of the polymer also adds significantly to the overall cost of the resulting asphalt/polymer mix. Thus, cost factors weigh in the ability to meet the above criteria for various asphalt mixes. In addition, at increasing levels of polymer concentration, the working viscosity of the asphalt mix becomes excessively great and separation of the asphalt and polymer may occur.\nOne result of the high viscosities experienced at increased polymer concentrations is that it makes emulsification of the asphalt difficult. As is known in the art and used herein, emulsification of asphalt refers to forming an emulsion of asphalt and water. Asphalt emulsions are desirable in many applications because the emulsion may be applied at lower temperatures than hot-mix asphalts because the water acts as a carrier for the asphalt particles.\nFor example, hot-mix asphalts, mixes of asphalt, aggregate, and a single polymer, commonly are applied at a temperature of 350° Fahrenheit (F.) to 450° F. (177° Centigrade (C) to 232° C.) to achieve the requisite plasticity for application. In comparison, an asphalt emulsion typically may be applied at 130° F. to 170° F. (54° C. to 77° C.) to achieve the same working characteristics. Emulsified asphalt products are generally used to reduce the release of environmentally-harmful volatile organic compounds normally associated with asphalts diluted with light carrier solvents such as diesel fuel, naphtha, and the like. Emulsification basically requires that the asphalt and any desired performance-enhancing additives be combined with an emulsifying agent in an emulsification mill along with about 20 to 40 percent by weight of water. However, high polymer loading in asphalt produces high viscosities and melting points, making emulsification of the polymer-asphalt composition difficult.\nThe bitumen/polymer compositions are prepared in practice at polymer contents range from about 3% to 6% by weight of bitumen depending on the nature and the molecular weight of the polymer and the quality of the bitumen. Gelling of the bitumen/polymer composition, which is observed fairly frequently during the preparation of the said composition or while it is stored, occurs as soon as the polymer content of this composition exceeds the above-mentioned threshold. It is thus difficult, in practice, to produce non-gellable bitumen/polymer compositions with a high polymer content, which are more economical to prepare and to transport than bitumen/polymer compositions with a lower polymer content.\nIn view of the above, bitumen compositions, which simultaneously meet the performance criteria required for road paving, and which are substantially free of volatile solvent would be advantageous. Additionally, viscosity stable bitumen compositions would be particularly advantageous. Further, a method for efficiently introducing the polymer into the bitumen composition would be desirable. In preparing the composition, significant mixing is needed to insure the uniform addition of both the polymer and any crosslinking agents. The crosslinking agents are usually added as a dry powder and mixed with the asphalt compositions.\nAs can be seen from the above, the art is replete with methods to improve the mixing of asphalt and polymer compositions. The needed elements for the commercial success of any such process include keeping the process as simple as possible, reducing the cost of the ingredients, and utilizing available asphalt cuts from a refinery without having to blend in more valuable fractions. In addition, the resulting asphalt composition must meet the above-mentioned governmental physical properties and environmental concerns. Thus, it is a target of the industry to reduce the cost of the polymers and crosslinking agents added to the asphalt without sacrificing any of the other elements.\nTABLE 1Properties of Various Asphalt GradesAASHTO M-226TESTAC 2.5AC 5AC 10AC 20AC 30AC 40Viscosity@140° F., poise250 ± 50500 ± 1001000 ± 2002000 ± 4003000 ± 6004000 ± 800(AASHTO T-202)Viscosity @275° F.; cSt,125175250300350400minimum (AASHTO T-201)Pen. @77° F.; minimum22014080605040AASHTO t-49)Flash Point, COC325350425450450450Minimum ° F.Ductility After TFOT10010075504025(AASHTO T-179)@77° F., 5 cm/min,minimumViscosity After TFOT10002000400080001200016000(AASHTO T-179)@140° F., poiseminimumTESTAR1000AR2000AR4000AR8000AR16000Viscosity@140° F., poise1000 + 2502000 + 500 4000 + 10008000 + 200016000 + 4000(AASHTO T-202)Viscosity@275° F., cSt,140200275400500minimum (AASHTO T-201)Pen. @77° F., minimum6540252020(AASHTO T-49)Percent of Original—40455052Pen. @77° F., minimumDuctility @77° F.,100100757575minimum, 5 cm/min\nMethods of preparing stable bitumen polymer compositions addressing these considerations have been disclosed in U.S. Pat. No. 6,180,697, the contents of which are incorporated herein as if set forth in full. The method steps therein disclose addition of the crosslinking agents containing elemental sulfur, zinc mercaptobenzothiazole, zinc oxide and dithiodimorpholine (as a replacement for sulfur). It is desirable to extend the teachings of U.S. Pat. No. 6,180,697 to include crosslinking agents having enhanced properties or to crosslinking agents that reduce or eliminate hydrogen sulfide emissions from the resultant composition."}
{"text": "Japanese Patent Application No. 2001-284594 is expressly incorporated by reference in its entirely.\n1. Field of the Invention\nThe present invention relates to a body upper structure for a vehicle. Particularly, it relates to a structure around a joint between a roof side-rail and a center pillar for a vehicle. In the vehicle, the roof side-rail forms one fore-and-aft directional framework member on one side of the roof, while the center pillar forms one up-and-down directional framework member on one side of the vehicle body.\n2. Description of Related Art\nJapanese Patent Application Laid-open No. 2000-313353 discloses a body upper structure for vehicle. In this structure, an engagement rib is formed on a connecting part between a center-pillar outer member and the roof side-rail. While, an engagement recess is formed on another connecting part between a center-pillar inner member and the roof side-rail. In assembling, by engaging and welding the engagement rib with the engagement recess, the roof side-rail and the center pillar are integrated into one body.\nNow, it is known that a downward load is applied on the roof side-rail when the vehicle has an xe2x80x9coffsetxe2x80x9d front collision or when the roof collapses. On the other hand, when the vehicle has a side collision, the center pillar is subjected to a load in a direction to thrust up the roof side-rail.\nThe above-mentioned body upper structure has the roof side-rail and the center pillar formed into one body; nevertheless there still remains a possibility that when the vehicle has the above collisions, the roof side-rail or the center pillar is locally and greatly deformed since the structure does not take measures to effectively transmit a collision load from the roof side-rail to the center pillar, and vice versa.\nUnder the above circumstance, it is an object of the present invention to provide a body upper structure for a vehicle, which is superior in transmitting a load from the roof side-rail to the center pillar, and vice versa, thereby reducing a local deformation in the roof side-rail or the center pillar at the vehicle collision or at the collapse of the roof.\nAccording to the invention, the above-mentioned object is accomplished by a body upper structure for a vehicle, comprising:\na roof side-rail which is formed to have a closed cross section and also provided with a vertical rib bridging between an upper wall and a lower wall both forming the roof side-rail; and\na center pillar united with the roof side-rail, the center pillar having a pillar outer member and a shoulder part formed on an upper end of the pillar outer member so as to fit a successive corner defined between lower and outer walls forming the roof side-rail thereby to support the roof side-rail;\nwherein the pillar outer member is provided, on an interior side thereof, with a plurality of vertical ribs at an interval in a fore-and-aft direction of the vehicle, the vertical ribs extending from a base of the shoulder part downward while respective upper ends of the vertical ribs are abutting on the lower wall of the roof side-rail; and\nwherein the vertical ribs are formed to project inwardly in a width direction of vehicle, up to a position where, in view along an up-and-down direction of vehicle, the upper ends of the vertical ribs substantially intersect with the vertical rib of the roof side-rail."}
{"text": "1. Field of the Invention\nThe present invention relates to a matrix type display apparatus such as a liquid crystal display (hereinafter referred to as an LCD), and particularly to a technique for the enlarged display of a screen.\n2. Description of the Related Art\nFIG. 2 is a structural diagram of a conventional LCD.\nThis LCD displays image signals supplied from a personal computer or the like (referred to below as a PC) on a color screen having a size of 1024 pixels horizontally by 768 pixels vertically. This LCD comprises a control circuit 10, a display signal circuit 20, a scan signal circuit 30, and a liquid crystal panel 40.\nRed, green, and blue color image signals formed from R, G, and B signals, clock signals CLK showing timings of a sampling of these R, G, and B signals, horizontal synchronizing signals HSYN, and vertical synchronizing signals VSYN are fed to the input side of the control circuit 10 from an unillustrated PC. On the basis of these signals, the control circuit 10 then generates start signals EI, clock signals CK, R, G, B data, and strobe signals STB and outputs these to the display signal circuit 20. The control circuit 10 also has the function of generating start signals ST and clock signals CP and outputting these to the scan signal circuit 30.\nThe display signal circuit 20 is provided with a 1024 stage shift register 21 which corresponds to the 1024 pixels in the horizontal direction, data latches 22 and 23, and a display drive section 24. The shift register 21 sequentially shifts and holds the start signals EI based on the clock signals CK and the held contents of each stage are fed to each stage of the data latch 22 as latch signals L1, L2, etc up to L1024. The data latch 22 holds the R, G, B data fed from the control circuit 10 based on the latch signals L1 to L1024.\nThe data latch 23 stores the R, G, B data of the 1024 horizontal pixels held by the data latch 22 at the strobe signal STB timing. The display drive section 24 generates display voltages S1, S2 etc up to S3072 which correspond to R, G, B data of the 1024 pixels stored in the data latch 23 and outputs these display voltages.\nThe scan signal circuit 30 is provided with a 768 stage shift register 31 and scan drive section 32 for sequentially scanning and displaying the 768 pixels in the vertical direction in horizontal line units. The shift register 31 holds the start signals ST in the initial stage register as scan signals and then sequentially shifts these scan signals in accordance with the timing of the clock signals CP. The contents at each stage of the shift register 31 are fed to the scan drive section 32 and are output as scan voltages G1, G2, etc up to G768.\nThe liquid crystal panel 40 is provided with 768 X electrodes X1, X2, etc up to X768 arranged at equal spacing in a line direction, and 1024 groups (note that each group is made up of 3 electrodes corresponding to R, G, and B) of Y electrodes Y1, Y2, etc up to Y3072 arranged at equal spacing in a column direction. Color pixel display is performed in line units in accordance with display voltages S1 to S3072 applied respectively to the Y electrodes Y1 to Y3072 at locations where an X electrode Xj to which a scan voltage Gj (wherein j is a number from 1 to 768) is applied intersects with a Y electrode Y1 to Y3072.\nIn this type of LCD, when R, G, B signals, clock signals CLK, horizontal synchronizing signals HSYN, and vertical synchronizing signals VSYN are supplied from a PC, horizontal cycle start signals EI and, thereafter, in synchronization with the clock signals CLK, R, G, B data is sequentially transmitted in single pixel units from the control circuit 10 to the display signal circuit 20. In the shift register 21, latch signals L1 to L1024 are generated in accordance with the clock signals CK. R, G, B data is further sequentially held in the data latch 22 in synchronization with the latch signals L1 to L1024.\nWhen the R, G, B data of the 1024 pixels of a single line is held in the data latch 22, a strobe signal STB is output from the control circuit 10. Accordingly, the R, G, B data of the single line in the data latch 22 is also stored in the data latch 23. In the display drive section 24, display voltages S1 to S3072 are generated based on the R, G, B data of the 1024 pixels stored in the data latch 23 and are output.\nMoreover, start signals ST for each vertical cycle and clock signals CP for shifting the start signals ST and generating scan signals are output from the control circuit 10 to the scan signal circuit 30. When the start signals ST are supplied, the output signal of the initial stage of the shift register 31 of the scan signal circuit 30 is set to a level xe2x80x9cHxe2x80x9d for designating display and output signals of subsequent stages are set to a level xe2x80x9cLxe2x80x9d for designating non-display. Output signals of each stage of the shift register 31 are then sequentially shifted one stage at a time towards the rear in synchronization with the clock signals CP which correspond to the horizontal cycle. The output signals of each stage of the shift register 31 are fed to the scan drive section 32 and scan voltages G1 to G768 are generated for each line in accordance with the display/non-display and output.\nThe display voltages S1 to S3072 which correspond to the R, G, B data of a single line output from the display drive section 24 are fed to the Y electrodes Y1 to Y3072 of the liquid crystal panel 40. Moreover, the scan voltages G1 to G768 output from the scan drive section 32 are fed to the X electrodes X1 to X768 of the liquid crystal panel 40. The timing of the strobe signals STB fed from the control circuit 10 to the display signal circuit 20 is substantially identical to the timing of the clock signals CP fed to the scan signal circuit 30. Therefore, a single line which corresponds to the X electrode Xj driven by the scan voltage Gj output from the scan drive section 32 is displayed through the display voltages S1 to S3072 based on the R, G, B data for the single line stored in the data latch 23. In addition, all other X electrodes other than the X electrode Xj are placed in a non-display state.\nA screen is displayed by the X electrodes X1 to X768 being sequentially driven from top to bottom by the scan voltages G1 to G768 sequentially output from the scan drive section 32.\nHowever, this type of conventional LCD has the following problems.\nFor example, when an image signal supplied from a PC is for displaying on a 640 horizontal pixel by 480 vertical pixel screen, no pixel data exists for displaying at the right hand side and bottom of a liquid crystal panel 40 having 1024 horizontal pixels by 768 vertical pixels. Therefore, the problem has existed that the displayed screen is small and the display position is offset to the top left.\nIn order to solve this problem, attempts have been made to provide a processing device and frame memory corresponding to the resolution of the liquid crystal panel 40 in the LCD. Image signals fed from the PC are interpolated by this processing device and converted into image data for a 1024 horizontal pixel by 768 vertical pixel screen and displayed.\nHowever, in this type of method, because a large amount of frame memory and a high speed processing device are necessary, the problem arises that costs are increased and the amount of power needed to run the processing device at high speed also increases.\nThe present invention solves the above problems in the conventional technology and provides a matrix type display apparatus such as an LCD which has a simple circuit structure and which is capable of enlarging the display of a screen.\nIn order to solve the above problems, the first aspect of the present invention is a matrix type display apparatus comprising: display means having M number of X electrodes (wherein M is a number greater than one) arranged in parallel, N number of Y electrodes (wherein N is a number greater than one) arranged so as to intersect the X electrodes, and display elements provided at each point of intersection of the X electrodes and Y electrodes, the display means performing a matrix type display at the points of intersection of the N number of Y electrodes and the X electrodes which are driven by scan voltage in accordance with a display voltage applied to each Y electrode; scan drive means for generating the scan voltage in order to sequentially drive the X electrodes; and display drive means for holding image data which corresponds to the X electrodes driven by the scan drive means, and for generating the display voltage based on the held image data and driving each of the Y electrodes, wherein the scan drive means and the display drive means have the following structure.\nNamely, the scan drive means is provided with M number of flip-flop circuits (referred to below as xe2x80x9cFF circuitsxe2x80x9d) and switches for altering the connections of the M number of FF circuits and, in normal display mode, the scan drive means forms shift registers having M number of stages by connecting the M number of FF circuits in series and also generates the scan voltages based on each output signal of the M number of FF circuits and, in enlarged display mode, the scan drive means forms shift registers having m number of stages (wherein m is a number less than M) by connecting a portion of the M number of FF circuits in parallel at a constant ratio and also generates the scan voltages based on each output signal of the M number of FF circuits.\nMoreover, the display drive means is provided with N number of FF circuits and switches for altering the connections of the N number of FF circuits and, in the normal display mode, the display drive means forms shift registers having N number of stages by connecting the N number of FF circuits in series and also holds the image data based on each output signal of the N number of FF circuits and generates the display voltages in accordance with the image data and, in the enlarged display mode, the display drive means forms shift registers having n number of stages (wherein n is a number less than N) by connecting a portion of the N number of FF circuits in parallel at the constant ratio and also holds the image data based on each output signal of the N number of FF circuits and generates the display voltages in accordance with the image data.\nThe second aspect of the present invention is the matrix type display apparatus according to the first aspect, wherein switches in the scan drive means and the display drive means are formed from analog switches (referred to below as xe2x80x9cSWxe2x80x9d) controlled by mode signals which designate the normal display mode or the enlarged display mode.\nAccording to the first and second aspects of the present invention, because the above described structure is employed, the following operation can be performed.\nIn normal display mode, the M number of FF circuits in the scan drive means are connected in series by switches such as SW or the like. Scan voltages are sequentially generated from each FF circuits and fed to the X electrodes of the display means. The N number of FF circuits in the display drive means are also connected in series by switches. Image data based on output signals sequentially output from each FF circuit is held, and display voltages are generated in accordance with the image data and fed to the Y electrodes of the display means.\nIn contrast, in enlarged display mode, the M number of FF circuits in the scan drive means are connected in parallel and in series at a constant ratio by switches. Scan voltages are generated based on output signals from each FF circuit and fed to the X electrodes of the display means. The N number of FF circuits in the display drive means are also connected in parallel and in series at a constant ratio by switches. Image data based on output signals from each FF circuit is held, and display voltages are generated in accordance with the image data and fed to the Y electrodes of the display means.\nThe third aspect of the present invention is the matrix type display apparatus according to the first and second aspects, wherein the scan drive means is provided with a plurality of integrated circuits used for shifting which are connected in a cascade connection and, in the enlarged display mode, the scan drive means simultaneously drives outputs of scan voltages for front or rear non-display regions based on selection signals which correspond to the integrated circuit connection position.\nThe fourth aspect of the present invention is the matrix type display apparatus according to the third aspect, wherein selection signals for a plurality of integrated circuits in the scan drive means are directly fed from ground voltage or power supply voltage supplying the integrated circuits.\nAccording to the third and fourth aspects of the present invention, the following operation is performed in the scan drive means in the matrix type display apparatus according to the first and second aspects of the present invention.\nSelection signals corresponding to the connection positions at which a plurality of integrated circuits used for shifting are connected in a cascade connection in the scan drive means are directly fed, for example, from a power supply and ground voltage identical to the integrated circuits. In enlarged display mode, the output of scan voltages for front and rear non-display regions is stopped based on the selection signals."}
{"text": "Methods for installing synthetic fields are well-known. As is well known, the installation method is primarily a manual process and typically differs from field to field. Most synthetic fields include a plurality of strips of synthetic turf that are placed or laid out on a substructure or sub-base, such as the ground, upon which the field is permanently installed. The strips together make up a completed synthetic field after they have been laid out, processed and then secured. Some or all of the strips of synthetic turf typically require markings thereon, such as yard line markings, numbers or the like. When the fields are installed, some of the markings, such as logos, yard line numbers, and hash marks are inserted into the field after it has been placed on the subbase at the location where the field is to be permanently located. As can be appreciated, the position of each individual strip with respect to the overall synthetic field is important. Similarly, the location of the markings on each of the individual strips is critical to the overall appearance of the field. If the markings are inserted in the field improperly, they can be out of alignment with respect to the other markings on the field or can be individually crooked, which can make the field look sloppy.\nDue to the nature of these fields, their construction, and the nature of the steps required for their installation, the typical manufacturing and installation steps suffer from a variety of disadvantages, including those discussed above. Also, the fabrication, manufacturing, and installation steps are extremely time consuming and labor intensive. Further, these steps can be extremely variable from field to field as well as costly due to the extensive labor commitment required."}
{"text": "In modern belt retractors, both a reversible belt tensioning unit and an irreversible belt tensioning unit are used for driving the belt reel in wind-up direction, when a vehicle condition is regarded as critical, in order to eliminate a slack in the safety belt and bring the vehicle occupant into an upright position. The reversible belt tensioning unit is configured as a pretensioner, which in the case that the critical vehicle situation does not lead to an accident returns the belt reel to its starting position. The irreversible belt tensioning unit, however, only is activated in actual accident situations, which subsequently require the replacement of the entire belt retractor.\nIn addition, modern belt retractors include a force limitation system which limits the force maximally exerted by the belt onto the body of the vehicle occupant, in that it permits a certain withdrawal of belt webbing from a defined force acting in the belt webbing. If the irreversible belt tensioning unit now is activated as long as the reversible belt tensioning unit still is coupled with the belt reel, the succeeding limitation of force not only requires to overcome the resistance provided by the force limitation system, but also to turn back the reversible belt tensioning unit against its direction of action. As a result higher forces than desired and provided by the design of the force limitation system are generated in the safety belt.\nThis invention creates a belt retractor, in which during a force limitation phase after a tensioning effected by the irreversible belt tensioning unit, the reversible belt tensioning unit reliably is separated from the belt reel, but is again operable after such force limitation phase."}
{"text": "GPR38 is a 7-transmembrane, G-protein coupled receptor, with high affinity for the peptide motilin [Feighner et al., Science 1999, 284, 2184], suggesting that endogenous motilin exerts all or most of its activity via this receptor.\nMotilin is a 22 amino acid peptide found in large amounts within endocrine-like cells of the gastrointestinal tract, and especially in the duodenum-jejunum areas. During fasting, the peptide is known to be associated with the onset of Phase III migrating complex activity within the stomach [Boivin et al., Dig. Dis. Sci. 1992, 37, 1562], suggesting a role in the mechanisms of prokinetic activity. Motilin is also released from the gut during feeding, sham feeding, gastric distension or by oral or intravenous nutrient application [Christofides et al., Gut 1979, 20, 102; Bormans et al., Scand. J. Gastroenterol. 1987, 22, 781], suggesting additional roles for this peptide in the modulation of motility patterns during feeding.\nIn animals or in man, motilin has long been known to increase gastrointestinal motility, and promote gastric emptying and intestinal propulsion in an anal direction, during both fasting and fed conditions. This activity is thought to be primarily due to a facilitation of at least the cholinergic excitatory function of the gut [Van Assche et al., Eur. J. Pharmacol. 1997, 337, 267], perhaps also involving the activation of the vagus nerve [Mathis & Malbert, Am. J. Physiol. 1998, 274, G80]. In addition, higher concentrations of motilin directly evoke a small contraction of the muscle [Van Assche et al., Eur. J. Pharmacol. 1997, 337, 267].\nThe antibiotic erythromycin was shown to mimic the gastrointestinal activity of motilin, in addition to its previously-described antibiotic properties [see Peeters, in Problems of the Gastrointestinal Tract in Anaesthesia Ed., Herbert M K et al. Springer-Verlag, Berlin, Heidelberg 1999, pp 39-51]. More recently, erythromycin has been shown to activate the GPR38 receptor, confirming its ability to mimic the function of motilin [Carreras et al., Analyt. Biochem. 2002, 300, 146]. In addition, the availability of this non-peptide motilin receptor agonist has allowed at least some clinical studies to be undertaken in order to examine the clinical potential of motilin receptor agonists. These studies have consistently demonstrated an ability to increase gastric emptying in various conditions associated with gastroparesis, such as functional dyspepsia and diabetic gastroparesis. Further, erythromycin has been shown to increase lower esophageal sphincter pressure in man, which together with the increase in gastric emptying, suggests a role in the treatment of gastroesophageal reflux disorders (GERD). Finally, erythromycin has been used to promote intestinal propulsive activity, finding clinical utility in the treatment of pseudo-obstruction and in conditions with impaired colonic motility [Peeters, in Problems of the Gastrointestinal Tract in Anaesthesia Ed., Herbert M K et al. Springer-Verlag, Berlin, Heidelberg 1999, pp 39-51].\nConsequently it is expected that agonists at the GPR38 receptor will mimic the activity of motilin or of other substances acting at this receptor, such as erythromycin, and find clinical utility in the treatment of gastrointestinal disorders associated with hypomotility, especially the functional bowel disorders such as GERD, functional dyspepsia (FD) and irritable bowel syndrome (IBS). The compounds will also be useful for the treatment of other GI conditions where the cause is known and in which GI motility is reduced. Such conditions include constipation, caused by various diseases such as those associated with neuropathy, and/or by the administration of other drugs, intestinal pseudo-obstruction, paralytic ileus following surgery or some other manipulation, gastric stasis or hypomotility caused by various diseases such as diabetes and/or by the administration of other drugs, or in enterally fed patients. Interestingly, the ability of motilin or erythromycin to activate the vagus nerve, the association of this nerve with changes in feeding behaviour [e.g. Furness et al., Auton. Neurosci. 2001, 92, 28] and the chromosomal location of GPR38 [based on Ensembl: 13q21.1 (58.46-59.46 Mb)] within the markers (D13S257-13q14.11 to D13S258 at 13q21.33) of a locus associated with obesity [Feitosa et al, Am. J. Hum. Genet. 2002, 70, 72] also suggests that agonists active at the GPR38 receptor will, in addition to promoting gastrointestinal motility, facilitate eating behaviours in at least those patients in which some degree of appetite suppression or cachexia is present. Such activity indicates that agonists at this receptor will find clinical utility in the treatment of symptoms associated with—for example—the treatment of cancer or by the presence of the cancer itself.\nIn addition to the ability of motilin receptor agonists to promote gastrointestinal motility, the association of motilin gene polymorphism with Crohn's disease [Annese et al., Dig. Dis. Sci. 1998, 43, 715-710] and the changes in motilin receptor density during colitis [Depoortere et al., Neurogastroenterol. Motil. 2001, 13, 55] suggests a utility for agonists at the motilin receptor for the treatment of inflammatory bowel conditions in general.\nFinally, GPR38 is also found in regions outside the gastrointestinal tract. These areas include the pituitary, adipose tissue, urinary bladder and certain areas of the brain. The former suggests clinical utility in the promotion of pituitary function, such as the release of growth hormone secretagogues, the presence within adipose tissue again suggests a role in the control of body weight, and the presence within the urinary bladder suggests a role for agonists at this receptor in the treatment of incontinence. The presence of GPR38 within the brain supports the gastrointestinal and feeding utilities already mentioned, but in addition, suggests an involvement of the receptor in a greater spectrum of vagal-hypothalamic functions.\nWO9410185, EP838469, WO9823629, DE19805822, and U.S. Pat. No. 6,165,985 claim erythromycin derivatives targeting GPR38 for use in disorders relating to gastrointestinal motility. WO9921846, WO0185694, WO0168620, WO0168621, and WO0168622 disclose a series of small molecule antagonists of the GPR38 receptor. JP07138284 and EP807639 disclose peptide agonists. JP09249620, WO02092592, WO05027637, US2005065156 and Li et al., (2004, Journal of Medicinal Chemistry, 47(7) p 1704-1708) disclose series of small molecule agonists. WO 05012331 and WO 05012332 disclose macrocyclic compounds which are agonists or antagonists of mammalian motilin or ghrelin receptors. WO 06127252 discloses erythromycin derivatives.\nWO07/007,018 describes compounds of formula (A), which have activity as agonists of the GPR38 receptor\n\nWO07/012,479 describes compounds of formula (B), which have activity as agonists of the GPR38 receptor\n\nA structurally novel class of compounds has now been found which provides agonists of the GPR38 receptor."}
{"text": "1. Field of the invention\nThis invention relates in general to wellhead equipment for oil and gas wells, and in particular to an emergency casing hanger seal.\n2. Description of the Prior Art\nIn a well of the type concerned herein, a wellhead housing has a landing profile or shoulder within its bore. When running casing conventionally, a casing hanger is installed on the upper end of the string of casing. The casing hanger lands on the landing shoulder in the bore of the wellhead housing.\nAfter cementing, a seal is positioned between the casing and the wellhead housing. The seal locates between machined surfaces on the wellhead housing and the casing hanger. A tubing hanger may be installed over the uppermost casing hanger. The tubing hanger is normally held by lockdown screws if the wellhead housing is located at the surface. The tubing hanger secures to the tubing which extends into the well. A tubing hanger seal locates between the tubing hanger and the wellhead housing.\nOccasionally, the casing will not proceed smoothly to the bottom of the well. When this occurs, the casing hanger will not be properly positioned to land in the wellhead housing. Generally, when this happens, the casing cannot be retrieved to the surface and becomes stuck. In that case, the casing must be cut above the landing shoulder after cementing. In the prior art, the casing is supported by slips in the wellhead housing.\nA problem exists in sealing against the casing stub, because the casing does not have a smooth machined surface for receiving a seal. The casing outer diameter has a high dimensional variation. The outer diameter may be slightly oval shaped. The surface of the casing may have many defects, such as rust, pock marks and tong marks. Various seals have proposed for sealing against the casing stub. However, improvements in locking the seal in the wellhead housing are desired.\nAnother desirable feature would be a means that would prevent the casing from moving upward due to temperature increase in the well during production. Movement of the casing could have damaging effects on the seal. Also, it would be desired to have a tubing hanger lockdown that did not employ lockdown screws, and would allow the seal to be removed without removing the lockdown."}
{"text": "This invention relates to a variable-speed pulley and clutch adapted for use in snowmobiles and other small vehicles. Such pulleys are connected to the engine shaft and may operate at high speeds which subject the rotating parts to high centrifugal forces. The parts must be designed and constructed to withstand high centrifugal forces, and the pulley must be smooth and reliable in its response to operator-controlled changes in the engine shaft speed and to variations in operating conditions. The present invention relates especially to pulleys of the type in which centrifugal weights between the movable pulley flange and a reaction cone connected to the fixed pulley flange force the flange together and provide torque drive between the cone and the movable flange.\nIn accordance with the present invention, the fixed pulley flange is mounted on a pulley hub which has a cylindrical primary bearing surface at the opposite side of the pulley groove from the flange and has a hub extension of reduced diameter extending therebeyond and providing an elongated outboard bearing surface. The movable pulley flange is carried by a sleeve hub slidably and rotatably mounted on the pulley hub and its extension by means of an inboard bearing surface engaged with the primary bearing and an outboard bearing spaced therefrom and engaged with the bearing surface of the hub extension. The two hubs are counterbored at the inner diameter of the sleeve hub to define a spring chamber about the hub extension which contains a spring that acts on the hubs to urge the flanges apart. To prevent metal-to-metal contact between the two pulley flanges and consequent wear from relative rotation of such flanges, the two hubs are formed with axially-facing thrust faces, radially inward of the primary bearing sleeve and radially outward of the spring, which come together to limit closing movement of the pulley flanges.\nThe movable pulley flange is urged toward the fixed flange by a selected number of symmetrically disposed plate-like centrifugal weights mounted between the movable flange and a reaction cone fixed to the outer end of the hub extension. The weights are desirably molded plastic disks disposed in radial channels on the face of either the reaction cone or preferably the movable pulley flange, and frictionally engage a surface of revolution on the other of such cone and flange. The limit of opening movement of the flanges is determined by the engagement of the retracted centrifugal weights with the cone and flange, and in consequence the opening-limit position is fixed by the position of the reaction cone which is fixed on the hub extension. The spring applies a preload to the centrifugal weights so as to produce an initial frictional drive between the driven reaction cone and the rotatably mounted movable flange. The movable flange thus has a limited throw, between the closing-limit position determined by engagement of the thrust faces on the hubs and the opening-limit position determined by the position of the reaction cone.\nSelective insertion of one or more spacers between the hub thrust faces will shift the closing limit position and thereby open the spacing between the flanges to suit wider belts. Corresponding selective insertion of spacers between the hub extension and the pulley hub, or otherwise in the connection between the reaction cone and the fixed pulley flange, will correspondingly shift the position of the opening limit of the movable flange, by moving the position of the reaction cone. Suitable selection of spacers will maintain the same or secure a different throw of the movable flange in response to action of the centrifugal weights.\nDesirably, the hub extension is connected to the main hub by a threaded joint at which axially-facing faces on the two parts come together and which permits the insertion of the spacers between such faces so as to axially adjust the extension (and the reaction cone it carries) relative to the main hub and its fixed pulley flange. This joint is desirably at the end of the spring chamber and the spacers extend under the end of the spring. This has the effect of maintaining the same effective spring length, and the same preload on the centrifugal weights.\nThe pulley is adaptable to various applications, as by changing the number or mass of centrifugal weights, by shifting the opening limit or the closing limit or both limits of movement of the movable flange, by varying the shape of the reaction cone, etc. In one illustrated modification, the reaction cone is formed with an annular groove which restrains the centrifugal weights until a desired engine speed is attained."}
{"text": "Large networks can include thousands of routers to transfer data among multiple neighboring networks, sub-networks and/or network devices. Each router typically includes multiple ports, each of which corresponds to one or more interfaces and subinterfaces. Configuration data is used to define characteristics to be associated with each router interface and subinterface including, for example, a speed at which the router interface will transfer data, a data encapsulation technique to be used by the router interface, one or more network devices to which the router interface is to transfer data, etc.\nThe figures are not to scale. Wherever possible, the same reference numbers will be used throughout the drawing(s) and accompanying written description to refer to the same or like parts."}
{"text": "A variety of industrial controllers (or industrial computers) are available for performing a wide variety of commercial and industrial processes including, for example, programmable logic controllers (PLCs). Generally, an industrial controller is a specialized computer having input/output (I/O) circuitry by which the controller is in communication with sensors and actuators on the machines of the industrial process. The industrial controller executes a control program to read inputs from sensors on the machines implementing the process, and based on the values of those inputs and control logic of the control program, produces outputs to actuators to control the process. Industrial controllers are frequently modular in construction, so that the number of input and output circuits and the processor power can be tailored to the demands of the particular process.\nIndustrial controllers can be implemented using centralized or distributed schemes. In many traditional centralized approaches, a single controller is employed for a variety of purposes, including production control, planning and scheduling and supply chain management. However, the use of such centralized control schemes has become more difficult and inefficient with respect to controlling many modern industrial systems and processes. Many modern industrial systems and processes not only involve high levels of complexity but also need to be highly flexible and dynamically reconfigurable in terms of their organization and manners of operation. Yet it is often difficult to program a single centralized controller so as to achieve effective control over such complex systems and to allow for such flexibility and dynamic reconfigurability of operation by such complex systems.\nWhile centralized control schemes face these limitations, over the past several years it has been recognized that distributed systems technologies, such as service oriented architectures (SOA) or multi-agent systems (MAS), offer powerful methods and techniques for implementing a new generation of flexible and reconfigurable industrial control solutions. In general, such distributed systems employ multiple independent or semi-independent intelligent decision making components (e.g., “intelligent agents,” simply “agents” or sometimes simply “services”) that are each capable of knowledge gathering, reasoning, information exchange, and cooperation with one another based on service matchmaking and negotiation. By way of information and knowledge sharing that occurs among these components, cooperation and coordination among the components can be achieved. At the same time, because the control responsibilities are spread among, and can be shifted among, the different control components, greater flexibility and dynamic reconfigurability is afforded.\nDespite the aforementioned advantages of such distributed control systems, it should be recognized that the control components of those systems are still tightly coupled (rather than loosely coupled) with one another in terms of the automated gathering and integration of data, information and knowledge. In particular, specific groups of control components are often configured (or even dedicated) for handling particular tasks associated with different portions or levels of an industrial process. For example, a given group of control components may deal exclusively with the manufacture of a given portion of an overall product being manufactured, deal exclusively with transportation of materials/components between different portions of an assembly line, or deal exclusively with the overall “enterprise management” of a factory while not dealing with more specialized aspects of factory operation.\nWhile the control components within a particular group may have well-developed communication protocols or formats by which the control components of that group can communicate with one another, communications among different groups often is more limited, is far from being open and, indeed, can require human intervention to take place. This is the case because, among other reasons, the distributed systems tend to be programmed with a focus on performing particular tasks rather than on interoperability. Notwithstanding these limitations associated with conventional distributed control systems, there remains a desire to achieve better integration of different portions or operational segments of industrial systems, including different levels of components associated with industrial systems (e.g., those associated with the shop floor level up to those associated with enterprise management).\nAccordingly, it would be advantageous if an open system that enabled and facilitated communication and cooperation among distributed control components and groups of such components at a larger scale with less (or without any) human intervention could be developed."}
{"text": "Phase locked loops (PLUs) are circuits well known in the communication arts for synchronizing a variable local oscillator with the phase and/or frequency of a component of a transmitted signal. Typically such circuitry includes a phase detector which is responsive to the transmitted signal and the output of a local oscillator for generating a phase error signal proportional to a difference between a component of the transmitted signal and the oscillator output. The phase error signal is subjected to smoothing and coupled to control the oscillation rate of the variable oscillator thus forming a self-regulating closed loop system. In order to enhance the operation of the PLL, some systems include a second loop which generates an error signal proportional to the difference between the frequency of the variable oscillator and the frequency of the component of the transmitted signal. The frequency error signal is added to the phase error signal for controlling the rate of the oscillator. Nominally, the variable oscillator will achieve the desired frequency before phase lock is achieved, at which time the frequency error signal is substantially zero and the PLL is controlled by only the phase error signal.\nAn example of a frequency/phase responsive PLL is described by Canfield et al., in U.S. Pat. No. 5,159,292 entitled ADAPTIVE PHASE LOCKED LOOP. In the Canfield system an unlock detector is used to a switch in the frequency error signal path which actively disconnects the frequency error term when the system approaches phase lock. The active disconnection of the frequency term precludes noise in the frequency error signal from causing jitter in the phase of the signal provided by the variable oscillator. In order to detect the frequency unlock condition of the loop, the Canfield system includes a quadrature phase detector that accumulates samples of in-phase and quadrature components of the input signal and the unlock condition is determined by counting the number of times the accumulated in-phase samples change polarity during an interval (e.g., a field) and comparing this count value to a threshold value. In order to prevent the combined frequency and phase error signal from excessively altering the oscillator frequency when unlocked, the combined error signal is applied to a limiter before being low pass filtered and applied to the oscillator."}
{"text": "This invention relates broadly to the field of consumer absorbent paper products and other washroom products, which includes, among other products, toilet tissue, paper towels and liquid products such as soap. More specifically, this invention involves an improved system and method for collecting data on individual and aggregate usage of washroom products in a manner that is not detectable to persons who are using the product.\nA great deal of technical and market research is invested by manufacturers of high quality absorbent paper products, such as Kimberly-Clark Corporation, in designing and manufacturing products that are consistent as possible with consumer needs and preferences.\nUnfortunately, for reasons both cultural and logistical, consumer habit information for certain products such as toilet tissue has been notoriously difficult to obtain. For example, data relating to factors as the total amount of paper used, the duration of time over which paper is used, the number of discrete pulls on the paper taken by a user and the amount of paper taken by a user per discrete pull would be very helpful for both marketing and engineering purposes. However, there has heretofore been no efficient way of compiling such data.\nWhen collecting this type of information, it is important that consumers not know their activities are being monitored, since this may change their behavior. In addition, some consumers might become apprehensive at the thought of being monitored in this manner.\nIt is clear there has existed a long and unfilled need for a system and method for obtaining information on consumer tissue preferences and habits that is accurate, efficient and substantially undetectable by consumers.\nAccordingly, it is an object of the invention to provide a system and method for obtaining information on consumer tissue preferences and habits that is accurate, efficient and substantially undetectable by consumers.\nIn order to achieve the above and other objects of the invention, a system for collecting data on toilet tissue use at a particular location, includes, according to a first aspect of the invention, a sensor for sensing one or more characteristics of toilet tissue use at a particular dispensing location; and a recorder, in communication with the sensor, for receiving and recording data from the sensor, whereby toilet tissue use at the dispensing location can be monitored and studied.\nAccording to a second aspect of the invention, a system for analyzing individual and aggregate usage of toilet tissue use at a particular toilet stall, includes a determining system for determining whether a person is present in a toilet stall; a sensor for sensing one or more characteristics of toilet tissue use at the toilet stall; an analyzer for analyzing data that is supplied by the sensor; and a recorder for recording data from at least one of the sensor and the analyzer, the determining system, the sensor, the analyzer and the recorder being configured and positioned so as to be substantially undetectable to a user, whereby toilet tissue use at the toilet stall can be accurately monitored and studied.\nAccording to a third aspect of the invention, a system for collecting data on the use of a roll form absorbent commercial paper product at a particular location includes a sensor for sensing one or more characteristics of paper use at a particular dispensing location; and a recorder, in communication with the sensor, for receiving and recording data from the sensor, whereby paper use at the dispensing location can be monitored and studied.\nAccording to a fourth aspect of the invention, a method for collecting data relevant to an individual\"\"s use of an absorbent roll type paper product in a communal area includes steps of: (a) determining when a person enters an area that is adjacent to a paper product dispenser; (b) monitoring the person\"\"s use of the paper product in such a manner that the person is not aware of such monitoring; and (c) recording data obtained from the monitoring.\nA still further aspect of the invention involves a system for monitoring product usage at a particular location. The system comprises a product dispenser in which a sensor device is positioned to detect removal of the product. A local transmission unit is in operative communication with the sensor. The local transmission unit generates a data signal of predetermined format representative of at least product usage, and wirelessly broadcasts the data signal. A data collection unit located remote from the dispenser is operative to receive the data signal and store usage information representative thereof.\nAdditional aspects of the invention are achieved by a system for monitoring product usage at a plurality of locations. The system comprises a plurality of product dispensers, each of which includes a sensor device positioned therein to detect removal of the product. A plurality of local transmission units are also provided, each associated with a respective dispenser to receive usage information from the respective sensor. The local transmission unit responsively generates a data signal of predetermined format. A data collection unit located remote from the dispenser is operative to receive data signals from all of the local transmitters and store usage information representative thereof.\nA roll web product dispenser constructed according to the invention comprises a dispenser housing having a roll support mechanism located therein. A sensor element is operative to rotate as the roll web product is withdrawn by a user. The dispenser also includes a detector positioned to sense rotation of the sensor element. A local transmission unit in operative communication with the detector is positioned in the dispenser housing.\nA liquid product dispenser constructed according to the invention comprises a dispenser housing having a liquid product reservoir located therein. An actuator mechanism is also provided for dispensing a quantity of product from the liquid product reservoir. A sensor is operative to detect that the quantity of product has been dispensed.\nA folded web product dispenser constructed according to the invention comprises a dispenser housing adapted to maintain the folded web product in a stack. A sensor is positioned to detect removal of individual product sheets from the stack.\nThese and various other advantages and features of novelty which characterize the invention are pointed out with particularity in the claims annexed hereto and forming a part hereof. However, for a better understanding of the invention, its advantages, and the objects obtained by its use, reference should be made to the drawings which form a further part hereof, and to the accompanying descriptive matter, in which there is illustrated and described a preferred embodiment of the invention."}
{"text": "Installing flooring can be a very time consuming and labor intensive process. This is especially true when the flooring being installed is a wood floor which has to be installed using individual boards which typically interlock with one another and are secured to a subfloor. Typically, the panels of wood floors have a groove extending longitudinally along a side of the board and groove extending longitudinally along the opposite side of the board so that the tongue of one board is accepted by the groove of an adjacent board.\nFurther, when installing a wood floor, it is undesirable to have any marks on the top surface of the floor from the component securing the floor to the subfloor. Thus, nailer devices are used to drive a nail or staple at an angle through the side of the floorboard so that the floorboard can be adequately secured to the subfloor without having any blemishes on the top surface of the floorboard. Again, this is a time consuming process where the nailer device must be placed along the edge of the floorboard so that the securing component can be driven through the floorboard at the desired angle.\nThere are devices which are used to move along a floorboard being secured to a subfloor which aligns the nailer device to the floorboard. Typically, nailers which are used with these devices and nailers which are used independently of the devices, have a pad which rests on the top of the floorboard being secured to the subfloor so that the nailer is properly aligned with the side of the floorboard. However, as the device moves along the floorboard, in some instances, the pad of the nailer may create friction which may cause the device to pivot and become misaligned with the floorboard being secured to the subfloor.\nIn addition, floorboards typically come in different heights. When the floorboards are being secured manually, the operator can adjust the nailer device accordingly by aligning the nailer with the floorboard in the desired manner. However, when the nailer device is automated, the nailer may not be properly aligned so that the fastener extends through the side of the board in the desired manner.\nTherefore, it is desirable to develop a nailer device which is motorized so that the device moves along the floorboard and aligns the floorboard in order to secure the floorboard to the subfloor. It is also desirable to develop a nailer device that is adjustable in height in which the pad of the nailer does not contact the floorboard so that the nailer device does not become misaligned with the floorboard, and can function properly with floorboards having different heights."}
{"text": "A photocurable filling restoration material called a composite resin is generally widely used for restoration of a tooth damaged by caries from an aesthetic viewpoint. Such composite resin (dental curable composition) is a material containing a polymerizable monomer (monomer), a filling material (filler), and a photopolymerization initiator as main components and having further added thereto a pigment, an additive, and the like as necessary. In this case, a radically polymerizable monomer having radical polymerizability is generally used as the polymerizable monomer. Therefore, polymerization shrinkage inevitably occurs upon curing of the dental curable composition, that is, a polymerization reaction of the radically polymerizable monomer.\nHence, when the dental curable composition is filled into a cavity of a tooth that requires restoration and then subjected to polymerization curing, a gap is liable to be formed between the cured dental curable composition and an inner wall of the cavity. When such gap is formed, plaque is liable to be accumulated in the gap, and hence secondary caries is liable to develop. Further, the cured dental curable composition is liable to be detached out of the cavity. Accordingly, it is preferred that the polymerization shrinkage of the composite resin upon curing be as small as possible.\nAs a method of suppressing the polymerization shrinkage of the dental curable composition upon curing, there are given: a method involving using a radically polymerizable monomer having a high molecular weight; and a method involving using a cationically polymerizable monomer having a smaller polymerization shrinkage than that of the radically polymerizable monomer and subjecting the polymerizable monomer to polymerization curing through cationic ring-opening polymerization. Further, in addition to the foregoing, there are also proposals concerning methods involving using dendritic polymers having dendritically branched molecular chains, such as a hyperbranched polymer and dendrimer (see Patent Literatures 1 to 3).\nThe dendritic polymers are broadly divided into, for example, a hyperbranched polymer and dendrimer. A conventional polymer generally has a linear shape, while such dendritic polymer has actively introduced branches. Therefore, as compared to the conventional polymer having a linear shape, the dendritic polymer has, for example, the following features: 1) having a unique molecular structure; 2) having a size on a nanometer order; 3) being capable of forming a surface retaining a large number of functional groups; 4) being capable of being reduced in viscosity as compared to the linear polymer; 5) exhibiting a fine particle-like behavior because of less entanglement of molecules; and 6) having voids in the molecule. In this case, the hyperbranched polymer has an advantage in terms of its ease of synthesis as compared to the dendrimer, and thus is particularly advantageous in industrial production. In general, the dendrimer is synthesized by repeating protection and deprotection, while the hyperbranched polymer can be synthesized in one stage and can be synthesized simply.\nFor example, in Patent Literature 1, there is a proposal of a technology involving using a (meth)acryloyl terminated hyperbranched polyester having at least one ethylenically unsaturated moiety as a polymerizable resin for forming a dental material for the purpose of, for example, reducing a polymerization shrinkage of the dental material. Further, in Patent Literature 2, there is a proposal of a technology involving using a dendritic polymer as a polymerizable composition that may be utilized as a dental material and is low in shrinkage. For example, in Patent Literature 2, there is a disclosure of an example using a hyperbranched polyesteramide as a hyperbranched polymer. In addition, in Patent Literature 3, there is a proposal of a technology involving using, as a dental material, a dendritic compound in which a core, first shell, and second shell for forming a molecule are bonded via a polyurethane group and the second shell is modified by a reaction with a (meth)acrylate. According to the patent literature, undesired polymerization shrinkage is advantageously affected by this technology.\nFurther, in addition to Patent Literatures 1 to 3, there is a proposal of a technology involving using, as a crystal component for forming a dental material, one having a dendritic, hyperbranched, or star-shaped structure (Patent Literature 4). According to the patent literature, the crystal component to be used for the dental material preferably has a hydroxyl group at the end."}
{"text": "1. Field of the Invention\nThe present invention relates to a chip mounter, and more particularly to a method for recognizing corners of an angular component such as a square component in the chip mounter to be mounted on the printed circuit board (PCB).\n2. Prior Art\nIn general, PCBs are boards on which several electronic components are integrated. The dense wiring of the electronic components on the PCBs can be accomplished through a printing method. Recently, single pattern PCBs are being miniaturized with high density double faces or multiple layers so that the single patten PCBs can be manufactured in large quantities at one time. Moreover, an industrial robot equipped with an apparatus which can perform a sight function that is highly adaptable and efficient has been proposed.\nA chip mounter includes an image recognition device constituted by a cognitive camera, an image signal processor, and a microprocessor. The chip mounter recognizes a posting of chip components which adhere to a nozzle and controls corners and an angle of the chip components at the next station according to the recognition result. In order to exactly recognize and correct the location and the angle of the square component 104 which adhere to the nozzle, the chip mounter controls the image recognition device to take a picture with respect to an image of the square component 104. The chip mounter judges a posting of chip components through the image thereof and corrects the chip component to an accurate posture by using the judged posting thereof. Accordingly, an operation for recognizing the posting of the chip components which are adherent to the nozzle should be performed. A performance of the chip mounter is influenced by how precisely and quickly the recognition process is performed. If the recognition operation is carried out in an abnormal state, the chip mounter cannot mount the chip components at the corresponding position by component data which are stored therein, resulting in a deterioration of the quality of the produced goods.\nU.S. Pat. No. 4,674,869 issued to Timothy R. Pryor et al. discloses a method and apparatus for electro-optically determining dimension, location, and altitude of objects. A first embodiment of the method and apparatus describes optical triangulation based coordinate measurement of complex surfaces, such as gear teeth and turbine blades. Another embodiment provides highly useful sensors for robot guidance and related purposes. However, the patent of Timothy R. Pryor et al. does not disclose a method for recognizing corners of a square component."}
{"text": "1. Field of the Invention\nThis invention relates generally to the production of stable cockroach aggregation pheromone products, active extracts thereof, and insect traps utilizing them. These products are useful in the control of cockroach infestation.\n2. Description of the Background\nBecause of their exposure to unsanitary areas, such as sewers, cockroaches are generally associated with the arousal of many health related problems. Over one hundred different species of disease-causing pathogenic bacteria have been found to be carried by cockroaches. Cockroaches have also been found to carry, in addition, other microorganisms such as viruses, protozoa, and parasitic worms (Cornwell, The Cockroach, vol. 1, (1968)).\nThe use of certain chemical and food attractants for cockroaches is known. Among them are starch, sugar, and fatty acids, such as myristic acid, and palmitic acid and esters of these fatty acids (Japanese Patent Publication No. 44-320). Other chemicals, such as derivatives of naphthalene (Japanese Patent Publication No. 61-69701 and 61-72702); alkylphenol or alkylphenylamine (U.S. Pat. No. 4,911,907), which are specific for male cockroaches, and furanose derivatives (U.S. Pat. No. 4,945,107) have been suggested. Most of the substances, however, have been shown to have low activity and had to be used in extremely high concentrations. In addition, a method for controlling cockroach infestation has also been proposed, which relies on exposing the cockroaches to irradiation with ultraviolet rays. For example, ultraviolet light of longer wavelengths and below 291 nm were utilized in an on and off manner for 15 min. each, for a total time of 9 to 96 hours. This method is said to also inhibit the formation of ootheca. However, this method of controlling cockroaches is expensive and difficult to implement except in specific instances.\nSeveral cockroach traps are now commercially available. Most of these traps, however, are rather ineffective because they do not contain effective cockroach attractants. Such is the case of sticky traps and traps containing boric acid bait.\nAggregation pheromones are produced by many cockroaches, such as the American cockroach or Periplaneta americana, the oriental cockroach or Blattela orientalis, and the German cockroach or Blattela germanica, among others. The production of these pheromones induces the cockroaches to aggregate locally (Bell et al, J. Kans. Entomol. Soc. 45: 414 (1972). The aggregation pheromone composition produced by the German cockroach has been reported to contain over 150 different chemical compounds (Fuchs et al, Z. Angew. Entomol. 99: 499 (1985)). Several components of the aggregation pheromone composition produced by the American cockroach have also been identified (Koehler, et al, Handbook of Pest Control, Franzak and Foster Co., Cleveland, Ohio (1990)). However, they were found to be very volatile and extremely expensive to produce. Thus, no method for producing large quantities of stable aggregation pheromones is currently unavailable.\nAggregation pheromones had previously been extracted from the feces of German cockroaches (Blattella germanica) using Freon 11 as a solvent (Fuchs et al., Z. Angew. Entomol. 99: 499 (1985)). The use of fecal materials from cockroaches not only poses logistical problems, but is also uneconomical. The Freon 11 extracts of cockroach feces have been characterized. They contain chemicals which can be roughly classified in two groups. The first group comprises volatile substances that attract cockroaches while the other group comprises non-volatile substances having arresting activity, e.g., act to aggregate the cockroaches.\nAn ideal attractant should not only attract cockroaches in low concentrations but, once cockroaches are attracted to the site, should have a prolonged stationary effect on the insects. However, up to the present time, even though aggregation pheromones were known to be present in the feces of cockroaches, the extraction of these pheromones has been impractical and uneconomical because it requires large amounts of fecal material and poses sanitary problems.\nU.S. Pat. No. 4,337,592 to Hasegawa discloses the extermination of cockroaches with ultraviolet radiation. In addition, the prior patent describes the elimination of aggregation pheromone activity by ultraviolet irradiation of insect feces and pheromones extracts from P. americana, P. fuliginosa, P. austrasiae, B. germanica, and N. cinerea. In the prior art patent, pheromones were obtained by extraction from cockroach excrement with methanol, and then irradiated with a holder type ultraviolet germicidal lamp.\nU.S. Pat. No. 4,608,774 to Sherman discloses a gravity operated cockroach trap having a rectangular funnel-like interior chute, with walls arranged at an angle greater than 10.degree. and preferably 44.degree. to 45.degree. so that the insect on the incline will be drawn by gravity into the receptacle. In one of the embodiments, the prior patent contemplates coating the chute's surfaces with silicone or vaseline. These are solid lubricants which are rendered less and less useful as the angle of the receptacle walls increases. When the inclination of the walls is much greater than 45.degree., the insects are able to scratch the surface of the solid coating and grip onto the solid wall material.\nThe complex nature of cockroach aggregation pheromone compositions, their instability, and an inability to mass-produce them in large quantities have precluded the manufacture and availability of effective cockroach lure or trap products containing them.\nThus, there exists a need for a low cost, reliable method of manufacturing formulations containing cockroach aggregation pheromones which preserve their ability to attract and arrest cockroaches."}
{"text": "1. Field of the Invention\nAbsorption refrigeration systems having means for automatically controlling the energizing medium.\n2. Description of the Prior Art\nIn U.S. Pat. No. 3,895,499 issued to N. E. Hopkins on July 22, 1975, there is described an absorption refrigeration system in which the cooling tower water is essentially uncontrolled and allowed to vary as influenced by the temperature of the ambient air flowing through the cooling tower in contact with the water spray devices. In the Hopkins system, the heat input to the generator is controlled in response to the entering or leaving chilled water temperature as is conventional in absorption machines. The basic concept embodied in the Hopkins system is a means for transferring solution from one point to another to prevent a dangerous operating condition.\nIn U.S. Pat. No. 4,164,128 issued to A. B. Newton on Aug. 14, 1979, there is described an absorption refrigeration system which is adapted to be powered by solar energy or a source of relatively low temperature waste heat. A control system automatically adjusts the maximum energizing temperature as the temperature of the cooling water and the temperature of chilled water supplied to the load vary from time to time.\nAn article by A. B. Newton entitled \"Optimizing Solar Cooling Systems\", ASHRAE Journal, November 1976, pp. 26-31, suggests that plural compartments for the storage of heated medium at different temperatures may be used in a solar powered absorption system. No specific means are described for controlling the selection of the storage compartments."}
{"text": "1. Field of the Invention\nThe present invention relates to powerline communications and in particular, powerline communication devices, and systems of use therefore.\n2. Description of the Related Art\nWith the growing need for the exchange of digital content (e.g. MP3 audio, MPEG4 video and digital photographs) there is a widely recognized need to improve digital communication systems. Powerline communication (PLC) is a technology that encodes data in a signal and transmits the signal on existing electricity powerlines in a band of frequencies that are not used for supplying electricity. Accordingly, PLC leverages the ubiquity of existing electricity networks to provide extensive network coverage. Furthermore, since PLC enables data to be accessed from conventional power-outlets, no new wiring needs to be installed in a building (or different parts of a building). Accordingly, PLC offers the additional advantage of reduced installation costs.\nReferring to FIG. 1, a household 100 typically has a distributed mains wiring system (not shown) consisting of one or more ring mains, several stubs and some distribution back to a junction box 112. In other constructs the distributed mains wiring system has a breaker box with circuits routed there from in a star configuration. For the sake of example, the household 100 has four rooms 104, 106, 108, and 120. Each room 104, 106, 108, and 120 may have a different number of outlets and other mains connections. For example, room 104 may have only one connection 122, room 106 may have two connections 124, 126, room 108 may have three connections 128, 130, 132 and room 120 may have six connections 134, 136, 138, 140, 142, 144.\nAccordingly, there are a variety of distances and paths between different power outlets in the household 100. In particular, the outlets most closely located to each other are those on multi-plug strips, and the outlets furthest away from each other are those on the ends of stubs of different ring mains (e.g. power outlets in the garden shed and the attic). Communication between these furthest outlets typically passes through the junction box 112. Nonetheless, the majority of outlets associated with a particular application (e.g. Home Cinema) are normally located relatively close together.\nBecause the channel capacity of a powerline and connectors attenuates according to, amongst other features, the frequency of a transmitted signal, current generation PLC systems have been developed to transmit signals at relatively low frequencies (i.e. below 30 MHz) and thereby obtain suitable transmission distances. However, the use of such low transmission frequencies limits the maximum data throughput obtainable by PLC systems. Only recently have powerline systems extended beyond 30 MHz, which causes problems due to regulations that require lower injected power above 30 MHz. This requirement places additional demands on the dynamic range of transceivers servicing PLC communications in these higher frequency ranges.\nThe processes of receiving analog signals and injecting modulated signals are treated differently by different PLC communication standards. A number of powerline communication standards have been defined. These include the Homeplug 1.0/1.1 standards, the Homeplug AV standard, the CEPCA standard, the Digital Home Standard, IEEE 1901, and ITU-T G.9960. Current PLC approaches often include some analog conditioning to the signal-path (e.g., low-pass filtering for anti-aliasing or smoothing, or AC coupling to remove the low-frequency [<<1 KHz] high voltage content of the electricity mains). However, because differing PLC communication standards support differing communication bands, differing modulations, channel bandwidths, etc., PLC devices typically service only a single PLC communication standard. In common with most communication systems, one of the main problems with prior art PLC systems is obtaining high throughput and wide coverage at reasonable implementation cost, whilst maintaining compatibility with existing technologies. There is, therefore, a need for improved PLC systems that overcome the above and other problems.\nCommunications within the household 100 or within other premises may also be serviced by a Wireless Local Area Network (WLAN), a cellular network, millimeter wave communications, e.g., 60 GHz, Wireless Personal Area Network (WPAN), Cable Modem Network, Local Area Network (LAN), and other communication techniques. Each of these communication types has its respective benefits and shortcomings. None of these communication types is typically able to provide a full coverage solution within the household 100 (or other premises). The shortcoming of all wired technologies is the lack of mobility thereof. Shortcomings of all wireless technologies are coverage holes, which are typical, interference from other wireless devices, including competing wireless devices, Radar, etc., and bandwidth limitations."}
{"text": "This invention is directed to an improved bathtub which includes an apron panel holding member located in a lip on the rim of the tub allowing an apron panel to be attached to either or both sides of the tub. Further, the tub is suspended from a supporting surface by one or more runners allowing the plumbing pipes associated with the tub to be located directly underneath a tub above a supporting surface.\nWith the advent of indoor plumbing, bathtubs were developed which included integrally formed drain holes allowing the tub to be conveniently drained. These tubs were generally supported on legs and were free-standing in their location. Further modifications of the tub yielded the built-in bathtub wherein the tub fit within a space or alcove within the bathroom. Because, however, the plumbing for draining the tub is located at one end or the other of the alcove it was necessary to supply both right and left hand tubs for proper location of the drain and the outwardly facing side in these built-in applications.\nGenerally most tubs are constructed of steel with enamel overlays. This resulted in a fixture which is quite heavy. With the recent boom in the mobile home market wherein weight is a factor to be considered in designing a mobile home, the use of extremely heavy fixtures is precluded. In a mobile home the plumbing cannot be located below the floor of the mobile home since this would place it on the bottom of the mobile home and therefore susceptible to road hazards when the mobile home is moved. This thus requires designers and builders of mobile homes to place the tub in a position in a mobile home dictated solely on where they can place the plumbing pipes and not dictated by design principles which allow for maximum utilization of living space and convenience of location of the tub for the owner of the mobile home.\nAside from the normal three-sided alcove placement of tubs many people prefer a sunken tub, a free-standing tub or possibly a tub only bordered against two or even just one wall. Each of these different locations of a tub presently require an individual tub which is solely designed only for that specific location, i.e., a sunken tub will be designed differently than a tub that is used against one wall, etc. This requires not only the manufacture of many different shapes and designs of tubs, but also requires a plumbing supply house to maintain a large inventory of diferent tubs each only utilizable in one particular construction. Because of the different shapes of these tubs shipping of these tubs is more costly because they cannot be orderly stacked and shipped within a small space but each almost requires its own separate packing container.\nIn view of the above discussion it is considered that there exists a need for an improved bathtub which can be utilized for a variety of installations including mobile home installations without regard to whether it will be placed against three walls, two walls, one wall, be a sunken tub or be a right or left handed tub. Further, there exists a need for a tub for mobile home installations which does not require the placement of a drain pipe below the structural support but which allows placement of the drain pipe between the structural support and the tub."}
{"text": "A conventional prior art hydraulic damper or shock absorber comprises a cylinder defining a working chamber having a piston slidably disposed in the working chamber with the piston separating the interior of the cylinder into an upper and a lower working chamber. A piston rod is connected to the piston and extends out of one end of the cylinder. A first valving system is incorporated for generating damping force during the extension stroke of the hydraulic damper and a second valving system is incorporated for generating damping force during the compression stroke of the hydraulic damper.\nVarious types of damping force generating devices have been developed to generate desired damping forces in relation to the speed and/or the displacement of the piston within the cylinder. These multi-force damping force generating devices have been developed to provide a relatively small or low damping force during the normal running of the vehicle and a relatively large or high damping force during maneuvers requiring extended suspension movements. The normal running of the vehicle is accompanied by small or fine vibrations of the un-sprung mass of the vehicle and thus the need for a soft ride or low damping characteristic of the suspension system to isolate the sprung mass from these small or fine vibrations. During a turning or braking maneuver, as an example, the sprung mass of the vehicle will attempt to undergo a relatively slow and/or large vibration which then requires a firm ride or high damping characteristics of the suspension system to support the sprung mass and provide stable handling characteristics to the vehicle. Thus, these multi-force damping force generating devices offer the advantage of a smooth steady state ride by eliminating the high frequency/small excitations from the sprung mass while still providing the necessary damping or firm ride for the suspension system during vehicle maneuvers causing larger excitations of the sprung mass.\nThe continued development of hydraulic dampers includes the development of multi-force damping force generating devices which are simpler to manufacture, can be manufactured at a lower cost and which improve the desired force generating characteristics."}
{"text": "The market for real estate information analysis is in its infancy. Traditionally, the real estate model has been broker-centric. The local real estate broker controls much of the localized information about a property becoming a necessary middle man for projects ranging from buying a house to selling a shopping center. With the advent of the Internet and its efficient method for disseminating information, more and more real estate information has become publicly available. This information, however, is spread out in multiple formats on various websites, databases, and other sources. This makes it very difficult, time consuming and expensive to compile sufficient information to make real estate purchase decisions. Real estate opportunities are often missed because of the time it takes to get actionable information on a site.\nTraditionally, real estate brokers and site selectors find potential commercial real estate sites for developers and tenants. One conventional method used by site selectors to find potential sites in a new area involves driving around the area and noting the location and quality of the different neighborhoods, and the location and quality of the existing commercial corridors. Thus the site selector will try to derive a potential store location's quality based on his observed quality of the surrounding area. The significant amount of time required to become familiar with an area is a reason why many developers and tenants turn to local real estate brokers for help. Another conventional method used by site selectors is to mark locations of existing retail stores on a map. This conventional method helps the site selector determine how far away a potential site is located relative to an existing retail store. One known approach used by commercial real estate site selectors makes use of paper maps with stickers to indicate store locations. Another known approach is to disseminate copies of books with hand marked store locations.\nOnce a developer or tenant's site selector finds a potential site, a conventional method involves ordering a demographic report from an in-house specialty team or outside consultant service. The demographic data for the potential site is then compared to a tenant's stated demographic requirements to determine demographic viability of a site. Additionally, tenants will often state how close they are willing to place stores together. This stated distance is compared to the distance the site selector marked for the site on his paper map to the known store locations. If the demographics and closest store distance meet the stated requirements for a particular retailer, then the developer will often move forward with plans to acquire and present the site to a tenant for development. These conventional methods must be repeated for each potential site, creating significant time and cost inefficiencies.\nThe problem with this conventional method of determining site viability is that there is a significant information and time gap between the site selector's first observation of a site and the developer's acquisition decision.\nIn view of the foregoing, there is a need to overcome the limitations of the conventional methods for finding site locations and determining site viability. There is a need to efficiently inform a site selector of the quality and location of the neighborhoods and commercial corridors without the site selector having to drive throughout the area or depend on a local broker. There is a further need to inform the site selector of a retailer's demographic and closest store distance requirements in a localized region, not just a generalized stated requirement. There is a further need to inform the site selector at the time of first observation whether the site's demographics and location to the nearest existing retail store meet the requirements of a particular retailer in a particular region. There is a further need to create a standard unified model to collect and disseminate site information throughout a commercial development organization to facilitate efficient site acquisition decisions."}
{"text": "1. Field of the Invention\nThe present invention relates to various power supply apparatuses used in automobiles, electric vehicles, construction machinery, various public welfare devices (such as video devices, television sets, and audio devices), various industrial devices (such as personal computers, communication devices, and FA control devices), and so on. Furthermore, the present invention relates to a power supply circuit and a large current load control apparatus of switching power supply using GaN-FETs.\n2. Description of the Related Art\nHeretofore, power semiconductor devices such as diodes, thyristors, triacs, GTO (Gate Turn Off) thyristors, bipolar transistors, MOS-FETs, and IGBTs (Insulated Gate Bipolar Transistors) are used in various power supply apparatuses. In these power semiconductor devices, a main current that flows through the power semiconductor device is controlled by switching control or analog control. These power semiconductor devices are devices serving as nuclei for implementing stabilized power supply apparatuses, such as switching regulators and linear regulators, and inverters for performing conversion to power having an arbitrary frequency and an arbitrary output voltage.\nIn these power semiconductor devices, there are switching loss caused by superposition of transient voltage and current at the time of switching and conduction loss caused at the time of conduction. These losses are converted mainly to heat. By the way, the conduction loss has such a characteristic as to become small as the on-resistance is decreased. This on-resistance corresponds to composite resistance of the channel resistance, bulk resistance, and so on existing within the semiconductor except contact resistance between electrodes and semiconductor layer interfaces in the power semiconductor device. Heat generated by the power semiconductor device causes a temperature rise of the power semiconductor device itself. Due to this temperature rise, the power semiconductor device operates at high temperature. Due to this high temperature operation, heat generation of the power semiconductor device is promoted. Such positive feedback is caused. As a result, thermal destruction of the power semiconductor device is caused by thermal runaway.\nUsually in the power supply apparatus, therefore, the power semiconductor device itself is provided with a radiation mechanism and in addition, with a radiator such as radiation fins for radiating heat generated by the power semiconductor device. Furthermore, a radiation fan is provided in order to improve the radiation effect in some cases. Furthermore, there is provided a fail safe mechanism that senses the temperature of the power semiconductor device and stops operation of the power semiconductor device when the temperature rises to such a value as to cause thermal runaway.\nHowever, the radiator is formed of a good thermal conductive material, such as aluminum, in order to provide the radiator with a function of heat sink as well. This results in a problem that the whole power supply apparatus becomes large in weight and capacity. Especially as for mobile power supply apparatuses for vehicles or portable power supply apparatuses, emergence of power supply apparatuses reduced in size and weight is demanded strongly.\nFor example, in a conventional power supply apparatus shown in FIG. 16, a radiator 302 having a large weight and a large capacity is needed. The power supply apparatus shown in FIG. 16 is a DC-DC converter power supply apparatus of a vehicle. A MOS-FET using a Si semiconductor material is incorporated in the power supply apparatus as a switching element. The apparatus main body 301 encloses all the elements forming the power supply apparatus. On the top of the apparatus main body 301, a radiator 302 formed of aluminum is provided. On a junction interface between the radiator 302 and the apparatus main body 301, a MOS-FET, which is not illustrated, stick to the radiator 302. Heat generated by the MOS-FET is absorbed by the radiator 302, and radiated by fins disposed on the top of the radiator 302. Because of installation of the radiator 302, the weight and volume of the whole power supply apparatus become excessively large.\nFurthermore, the radiator needs to stick to the power semiconductor device in order to favorably transfer the heat from the power semiconductor device. This brings about limitation on design that the radiator needs to be disposed with due regard to the outer periphery of the casing of the power supply apparatus and the radiation path. This results in a problem that the degree of freedom of the power supply apparatus is reduced. In addition, as for devices such as vehicles using the power supply apparatus, design of the whole device must be changed according to the disposition position of the power supply apparatus. Thus, there is also a problem that the design of the whole device is largely affected.\nIn addition, in designing the radiator, it is necessary to conduct sufficient radiation design with due regard to the ambient environment of the power supply apparatus. In addition, it is necessary to prevent the power semiconductor device, which is a heat source, from affecting other circuit elements having low heat-resisting property. This results in a problem that much time and labor are required for radiation design and arrangement design of other circuit elements included in the power supply apparatus.\nFurthermore, a heat protection circuit for preventing thermal runaway of the power supply apparatus is needed. This heat protection circuit monitors the temperature changes of important components, such as the power semiconductor device, included in the power supply apparatus. When the temperature has risen to a predetermined value, the heat protection circuit conducts fail safe control, such as stopping the power supply apparatus and causing shift to a low dissipation mode. This heat protection circuit is a complicated circuit that senses the temperature, output current, and so on and conducts a shift to fail safe control by using a logic processing circuit. Thus, there is a problem that the power supply apparatus must have a heat protection circuit having such a complicated circuit.\nRecently as semiconductor devices having high heat-resisting property, high breakdown voltage, high operation rate and low conduction loss, GaN (gallium nitride) FETs (Field Effect Transistors) have been developed.\nHeretofore, such power supply circuits have been applied to, for example, automobiles, various public welfare devices (such as video, television and audio devices), and industrial devices (such as personal computers, communication devices, and FA control devices).\nThe above described power supply circuit includes a transformer. A transistor made of, for example, a power MOS element turns on and off according to a gate signal. As a result, an output voltage is generated on a secondary winding side.\nIn the above described power supply circuit, however, the power MOS element used as the transistor, such as a power MOS-FET (2SK2313) generates much heat. Therefore, it is necessary to perform the radiation design accurately. A channel temperature Tch-max of the power MOS-FET itself at an ambient temperature of 85° C. is calculated as\n                              T          ch_max                =                              T            a_max                    +                                    P              total                        ×                          R                              th                ⁡                                  (                                      ch                    -                    a                                    )                                                                                            =                  85          ⁢          °          ⁢                                          ⁢                      C            .                                                  ⁢                          +                                                          ⁢              2                                ⁢                                          ⁢          W          ×          50          ⁢          °          ⁢                                          ⁢                      C            .                          /                                ⁢          W                                        =                  185          ⁢          °          ⁢                                          ⁢                      C            .                              where Ta—max: ambient temperature\nPtotal: total loss\nRth(ch-a): thermal resistance between channel and environment.\nThe temperature rises up to the channel temperature or higher. Therefore, it is necessary to provide a radiation plate. Supposing derating of 50° C. for a channel temperature of 150° C., the radiation plate design is represented as\n                                          θ            f                    <                                    θ                              ch                ⁢                                  -                                ⁢                a                                      -                          (                                                θ                  i                                +                                  (                                                            θ                      c                                        +                                          θ                      s                                                        )                                            )                                      =                              7.5            ⁢            °            ⁢                                                  ⁢                          C              .                              /                                      ⁢            W                    -                      (                                          0.833                ⁢                °                ⁢                                                                  ⁢                                  C                  .                                      /                                                  ⁢                W                            +                              0.8                ⁢                °                ⁢                                                                  ⁢                                  C                  .                                      /                                                  ⁢                W                                      )                                                  =                  5.9          ⁢          °          ⁢                                          ⁢                      C            .                          /                                ⁢          W                    where θf: thermal resistance of radiator\nθch-a: total thermal resistance between channel and environment\nθi: thermal resistance between junction portion and case (internal thermal resistance)\nθc+θs: thermal resistance between case and radiator From the foregoing description, it is necessary to select a radiator having a thermal resistance of 5.9° C./W or less. For example, therefore, a radiation plate made of an aluminum plate of 100 cm2 having a thickness of 1 mm becomes necessary. As a result, the conventional power supply circuit has a problem that the circuit configuration becomes large and heavy because of the radiation plate.\nFurthermore, heretofore, such a large current load control apparatus is applied to, for example, lighting control of head lamps of automobiles.\nIn the above described large current load control apparatus, lighting control of a head lamp is conducted by turning on and off a power MOS-FET formed of, for example, an on/off control switching element provided on a power supply line, which connects a battery to the lamp, under the control of a microcomputer.\nIn this control apparatus, however, the power MOS-FET used as the switching element of on/off control generates much heat. Therefore, it is necessary to perform the radiation design accurately. A channel temperature Tch—max of the power MOS-FET is calculated as\n                              T          ch_max                =                ⁢                              (                          T              a_max                        )                    +                                    (                              R                on_max                            )                        ×                          (                              l                o_max                            )                        ×                          (                              l                o_max                            )                        ×                          R                              th                ⁡                                  (                                      ch                    -                    a                                    )                                                                                            =                ⁢                  85          ⁢          °          ⁢                                          ⁢                      C            .                                                  ⁢                          +                                                          ⁢              0.013                                ⁢                                          ⁢          Ω          ×          10          ⁢                                          ⁢          A          ×          10          ⁢                                          ⁢          A          ×          50          ⁢          °          ⁢                                          ⁢                      C            .                          /                                ⁢          W                                        =                ⁢                  150          ⁢          °          ⁢                                          ⁢                      C            .                              where Ta—max: ambient temperature\nRon—max: on-resistance\nlo—max: current value\nRth(ch-a): thermal resistance between channel and environment.\nThe temperature rises up to the channel temperature. Therefore, it is necessary to provide a radiation plate. Supposing derating of 50° C. for a channel temperature of 150° C., the radiation plate design is represented as\n                                          θ            f                    <                                    θ                              j                -                a                                      -                          (                                                θ                  i                                +                                  (                                                            θ                      c                                        +                                          θ                      s                                                        )                                            )                                      =                ⁢                              11.5            ⁢            °            ⁢                                                  ⁢                          C              .                              /                                      ⁢            W                    -                      (                                          0.833                ⁢                °                ⁢                                                                  ⁢                                                      C                                    .                                      /                                                  ⁢                W                            +                              0.8                ⁢                °                ⁢                                                                  ⁢                                  C                  .                                      /                                                  ⁢                W                                      )                                                  =                ⁢                  9.9          ⁢          °          ⁢                                          ⁢                      C            .                          /                                ⁢          W                    where θf: thermal resistance of radiator\nθj-a: total thermal resistance between channel junction portion and the outside air\nθi: thermal resistance between junction portion and case (internal thermal resistance)\nθc+θs: thermal resistance between case and radiator\nFrom the foregoing description, it is necessary to select a radiator having a thermal resistance of 9.9° C./W or less. For example, therefore, a radiation plate made of an aluminum plate of 6 cm2 having a thickness of 1 mm and a weight of approximately 10 g becomes necessary. As a result, the conventional large load control apparatus has a problem that the circuit configuration becomes large and heavy because of the radiation plate.\nTherefore, it is one object of the present invention is to provide a power supply apparatus capable of implementing reduction of the size and weight, conducting flexibly design including the radiation design, and remarkably reducing the time and labor required for the design.\nFurthermore, it is an another object of the present invention is to provide a power supply circuit capable of reducing the heat generated by the transistor, thereby making the radiation plate unnecessary, and implementing reduction of the size and weight of the circuit.\nFurthermore, it is still another object of the present invention is to provide a large current load control apparatus capable of reducing the heat generated by the on/off control switching element, thereby making the radiation plate unnecessary, and implementing reduction of the size and weight of the circuit."}
{"text": "1. Field\nEmbodiments of the disclosure relate generally to the field of control surface actuation systems and more particularly to embodiments for a shape memory alloy actuator with dual collinear shape memory alloy tubes having antagonistic reaction about a control temperature point.\n2. Background\nWind tunnel models typically require movable control surfaces to allow simulation of various control aspects of the vehicle being modeled. Unmotorized surfaces are often used due to their simplicity. However, such surfaces must be positioned by hand requiring interruption of testing to position the surfaces at desired control angles. Models are typically of reduced scale and therefore full size actuators which would be employed in actual vehicles are not readily adaptable for use. Various actuation systems have been employed in wind tunnel models including electromechanical actuators and shape memory alloy (SMA) actuators using wires for hinge moment actuation using differential pull from SMA wires. However, electromechanical actuation is relatively bulky because of low power densities and the need for complex electric motor/gear assemblies. As such, the amount of space required in the supporting structure (for example in a vertical tail) may limit the amount of instrumentation such as pressure sensors that can be installed in the model and may reduce the structural strength which tends to limit their use to lower pressure tunnels having lower loads. Lower pressure tunnels do not match the aerodynamic characteristics of a full scale airplane as well which limits their fidelity as design tools for testing aircraft configurations. SMA wire actuation has limited power and strength, and therefore is similarly suitable for low pressure wind tunnel testing only.\nIt is therefore desirable to provide actuators for use in models, or other space constrained applications, to improve wind tunnel test efficiency by reducing the number of times the wind tunnel has to be opened to complete model changes while providing an actuator with sufficient force capability for higher pressure wind tunnels. It is also desirable to provide an actuator having power density much greater than traditional solutions for actuated control surfaces allowing its implementation into applications requiring higher forces and into more restrictive spaces. It is further desirable to provide an actuator with the ability to place the actuating elements on the hinge line of a control surface."}
{"text": "In facilities for production of metal sheet materials, in particular steel-sheet-making process lines, when making a conveying roll rotate at a high speed to feed steel sheet, the phenomena of slipping and snake motion of the steel sheet, deposition and buildup of dirt on the surface of the conveying roll etc. occur.\nIn particular, in a continuous annealing furnace, the hearth rolls convey steel sheet in a high temperature state, so buildup easily occurs on the hearth roll surfaces. When such buildup occurs, the shapes of the buildup are transferred to the steel sheet surface whereby the surface quality is impaired and the grade of the steel sheet deteriorates. Not only this, at the time of periodic repair, it is necessary to clean the hearth rolls to remove foreign matter deposited on their surfaces, so this becomes one cause of a drop in productivity.\nThe buildup on the hearth roll surfaces is the phenomenon of the iron, manganese oxide, etc. on the steel sheet surface sticking to and building up on the hearth roll surfaces. To prevent this, it is effective to suppress the reaction of the sources of buildup, that is, iron, manganese oxide, etc. with the hearth roll surfaces or facilitate removal of reaction products.\nAs a measure for suppressing buildup on a hearth roll surface, a coating obtained by providing a flame sprayed alloy layer comprised of only a heat resistant alloy on a hearth roll, flame spraying particles of a carbide or a mixture of a carbide and oxide on the sprayed alloy layer, and depositing metal oxides comprised of Cr2O3 and Al2O3 on the outermost layer of the surface to give a chemically converted layer has been proposed (for example, see Japanese Patent Publication (B2) No. 8-19535).\nHowever, the Cr2O3 and Al2O3 provided at the outermost layer of the surface easily react with manganese oxide, so the coating has the problem of easy generation of buildup due to manganese oxide.\nFurther, a coating containing, by wt %, chrome carbides: 10 to 25% and Ni: 5 to 15% and having a balance of one or more of a carbide and boride of tungsten and unavoidable impurities has been proposed (for example, see Japanese Patent Publication (A) No. 3-86306).\nHowever, with this coating, WC oxidizes at a high temperature and the coating peels off, so there is the problem that use for a long time in an annealing furnace is not possible.\nFurthermore, a coating comprising a cermet material of a particle structure comprised of, by wt %, 50 to 90% of chrome carbides and a balance of unavoidable impurities and a nickel-chrome alloy, at least 70% of the carbide particles enclosed by the alloy, and having an average particle size of 5 to 100 μm has been proposed (for example, see Japanese Patent Publication (A) No. 6-116703).\nFurther, a cermet coating containing at least one of CrB2, ZrB2, WB, TiB2, and other borides in 1 to 60 vol %, containing at least one of Cr3C2, TaC, WC, ZrC, TiC, NbC, and other carbides in 5 to 50 vol %, and having a balance of substantially metal is provided (for example, see Japanese Patent Publication (A) No. 7-11420).\nHowever, these coatings contain chrome carbides resistant to buildup at a high temperature, but if used for a long time in a continuous annealing furnace, the metal ingredients and chrome carbides in the coatings react resulting in the coatings become brittle and the coatings peeling off.\nThat is, in the past, various coatings have been proposed, but none of these coatings have been able to completely prevent buildup on hearth roll surfaces."}
{"text": "Providing defense mechanisms against computer hackers and attacks on computing infrastructure has been an increasingly significant concern with the growing importance of network-based computing. With the growth of network-based computing over the Internet and the growing reliance on cloud computing services, protecting computers and servers connected to a network from malicious network-based attacks has become a significant issue for businesses and organizations. For example, hackers may attack a computing service over a network using denial-of-service attacks. SYN flooding is a particularly popular denial-of-service tactic. In a SYN flood attack, an attacker sends a succession of SYN requests to a target system in an attempt to overwhelm the target system's resources so that it becomes unresponsive to legitimate requests. Computing service providers must account for and protect against such attacks while maintaining an effective level of service to their customers."}
{"text": "1. Field of the Invention\nThis invention'relates to a method and apparatus for optimally recording or transmitting geographical position and events. In particular the present invention provides apparatus to determine position, time and speed of an object such as a vehicle and log data in accordance with changes in pre-determined parameters.\n2. Description of the Prior Art\nRecording or transmitting the position (latitude, longitude and possibly altitude) and events that occur for an object is a useful exercise in many scenarios in business. There are many examples of systems that today record and or transmit positional data, but they all suffer from the same problems. They do not correctly choose the positions that are to be stored or transmitted, so as to maximize the level of useful detail on a map but at the same time minimize the number of actual positions used to do this. Many systems currently base their logging on an interval (eg. every 60 seconds) or a distance (eg. every 200 m). When looking at this data on a map, very often these points don't show any new useful additional information or, which is worse they miss out a significant event like a turn or a big drop in speed.\nThe present invention is used to optimize exactly when these pieces of information are recorded or transmitted. This then minimizes the memory requirements if these points are stored, or it minimizes time or cost if the information is transmitted. In addition, the usefulness of the information is maximized by making sure that any significant events are caught, like turning a corner or a large change in speed."}
{"text": "A combined power generation system based on geothermal energy and solar energy in the prior art employs carbon dioxide as a working fluid and improves thermal efficiency of the system by extraction, recompression and regeneration, because the specific heat of carbon dioxide varies greatly with pressure. For example, within 200 degrees centigrade, the specific heat of carbon dioxide at a pressure of 15 MPa is greatly different from that of carbon dioxide at a pressure of 8 MPa. Consequently, the temperature difference of heat transfer in regenerator increases and thus the irreversibility of heat transfer of the system increases, since the heat capacity (a product of mass flow rate and specific heat) of the high-pressure and low-temperature carbon dioxide is different from that of the low-pressure and high-temperature carbon dioxide in a simple regenerative cycle. The heat capacity at the high-pressure side may be reduced to match the heat capacity of the carbon dioxide at the low-pressure side by recompression and partial cooling. The temperature difference of heat transfer can decrease in the regenerator, so as to make full recover of the heat energy in the exhaust gas at a turbine outlet.\nHowever, a compressor operating efficiently under high-temperature and high-pressure is needed to ensure the recompression and partial cooling cycle implement. Specifically, the working circumstance of the compressor used for supercritical CO2 recompression Brayton cycleis adverse. Although the pressure ratio is not great, the ratio of specific heat of carbon dioxide is large, and the carbon dioxide at the final state of recompression may have a temperature up to 200 degrees centigrade and an operation pressure over 20 MPa, which will cause many problems for compressor designing and manufacturing.\nCurrently, the compressor operating efficiently under high-temperature and high-pressure has not been commercialized yet, so the method for improving thermal efficiency of the existing combined power generation system based on geothermal energy and solar energy by recompression and partial cooling still remains in the test program. Even though the compressor operating efficiently under high-temperature and high-pressure can be created in future, the manufacture difficulty, the manufacture cost and the operation cost of the combined power generation system based on geothermal energy and solar energy will be increased considerably."}
{"text": "Vehicles are often provided with a plurality of seats for supporting users. Conventional seats are often provided with an ability to adjust a horizontal position of the seat to accommodate users of various heights and various comfort preferences. However, additional solutions are needed for seats in vehicles that enable greater degrees of customization to the user."}
{"text": "The present invention relates to hot water heaters, and relates more particularly to such a hot water heater which has heat exchanging devices connected in series to recycle heat energy for pre-heating water permitting water to be quickly heated to the desired temperature.\nConventional hot water heaters commonly comprise a single heat exchanging device to carry water, and a gas burner to heat the flat bottom of the heat exchanging device. During the operation of a hot water heater, hot air from the heat exchanging device is directly drawn away from the housing of the hot water heater into the air, causing the air polluted. Because cold water is not pre-heated before it is guided into the heat exchanging device, it takes a lot of time and consumes a lot of energy to heat water to the desired temperature. Furthermore, when the heat source is operated to heat the heat exchanging device, heat is transmitted to water from the bottom, therefore heat cannot be evenly distributed through the inside of the hot water heater. Because different parts of the hot water heater bear different temperatures, their structural strength will be differently affected, causing the service life of the hot water heater shortened."}
{"text": "A hard-disk drive (HDD) is a non-volatile storage device that is housed in a protective enclosure and stores digitally encoded data on one or more circular disk having magnetic surfaces. When an HDD is in operation, each magnetic-recording disk is rapidly rotated by a spindle system. Data is read from and written to a magnetic-recording disk using a read/write head that is positioned over a specific location of a disk by an actuator. A read/write head uses a magnetic field to read data from and write data to the surface of a magnetic-recording disk. A write head makes use of the electricity flowing through a coil, which produces a magnetic field. Electrical pulses are sent to the write head, with different patterns of positive and negative currents. The current in the coil of the head induces a magnetic field across the gap between the head and the magnetic disk, which in turn magnetizes a small area on the recording medium.\nMany HDDs are configured to include a ramp (also referred to as a load/unload ramp). Ramp load/unload technology involves a mechanism that moves the head stack assembly (HSA), including the sliders, away from and off the disks and safely positions them onto a cam-like structure. The cam typically includes a shallow ramp on the side closest to the disk. During a power-on sequence, for example, the read/write heads are loaded by moving the sliders off the ramp and over the disk surfaces when the disks reach the appropriate rotational speed. Thus, the terminology used is that the sliders or HSA are “loaded” to or over the disk (i.e., off the ramp) into an operational position, and “unloaded” from the disk (i.e., “parked” onto the ramp) such as in an idle position.\nFurthermore, many HDDs are configured with a latch mechanism to secure the HSA, or at least the suspension/slider portion of the HSA, on the ramp (or perhaps on a landing zone of the disk for HDD's absent a ramp) when not operating. Thus, the head slider is prevented from dislodging from the ramp and moving over the disk when the actuator rotates due to external vibration or shock events. In light of the customer demand and industry move toward thinner and thinner HDD's, internal space constraints present challenges for sub-system design and development.\nAny approaches described in this section are approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section."}
{"text": "1. Field of the Invention\nThe present invention relates to electronic circuits and systems. More specifically, the present invention relates to electronic sensor systems.\n2. Description of the Related Art\nFor systems having unknown delay variables in the control and/or data paths, data-capture clocks are often used to compensate for the variable delays. For instance, in a video imaging system having a detector and a signal processor, the signal coming into the processor from the detector may not be aligned with the system clock of the processor. A data-capture clock, which is typically a phase shifted version of the signal-processing clock, is used to capture the detector data at its optimum sampling time. However, this can result in data path uncertainty if the sampling edge of the data-capture clock causes the captured data to change during the setup and hold time region of the signal-processing clock. The relative position of the data-capture clock also causes undesirable timing shifts to the signal processing module. These shortcomings limit the adjustment range of the data-capture clock, which, for some applications (especially high speed applications), is critical to the system.\nMoreover, signal processing modules in a digital system are often synchronized to a common system clock. This implies that the optimum data-sampling time of all data entry paths all must fall within a certain region of the signal-processing clock. Costly custom tailored add-ons (i.e. delay lines) for an individual system may be needed to guarantee this timing condition.\nOne prior art solution to these data path uncertainty and timing shift issues includes separating the signal processing module into control and processing sections, each running its own tailored timing. This assumes that there cannot be much correlation between the two sections within the module because of possible clock alignment in the forbidden zone. For multi-module systems, it is hard to pass data among different sections because the processing sections in each module may have a different clock phase.\nA second prior art solution delays the detector controls to align the detector video with the signal-processing clock. This approach requires expensive, power hungry, and space consuming variable delay lines to obtain the desired delay adjustment resolution.\nAnother solution includes adding an intermediate data buffer stage to allow independent access to the data at pixel level timing. This requires large memory buffers, and the signal processing may not start at a dedicated clock cycle due to clock phasing issues. For systems having multiple parallel signal processing modules, signal processing can only begin when all the buffers are not empty because the clock phasing may not be the same for all modules.\nAnother prior art approach involves supplying the detector with fixed controls and having the detector electronics provide the data clock associated with its outputs. Sections in the signal processing module must run its own tailored timing. Data exchanges and communications among sections may be difficult, especially for multi-module systems, due to clock phasing issues.\nHence, there is a need in the art for an improved system or method for aligning an input data signal with the system clock of a signal processor, which overcomes the shortcomings of prior systems."}
{"text": "1. Field of the Invention\nThe present invention relates to a system for processing input information or controlling execution of a process in accordance with the environment around the system and/or the attributes of user and a method therefor.\n2. Related Background Art\nHitherto, a system has been known which interprets input in a natural language to infer a goal or an intention of a user from the input. To realize a system of the foregoing type, the system must have a natural language parser (analyzer) for converting a natural language, such as English, into a computer expression and a mechanism, which uses a knowledge base formed by making a goal and a plan for achieving the goal to correspond to each other to infer a goal behind input.\nAnother system is available in which input in a natural language cannot be processed but which performs a predetermined process in accordance with commands, a menu or the like.\nIn the conventional natural language processing system, input in a natural language is converted into a computer expression and the goal behind that input is inferred using a knowledge base which has the knowledge of goals and plans. However, it has not been applied to predict the results of a possible plan, especially its undesirable effects.\nIn a system employing an input method which does not use the natural language, a command is interpreted to a limited meaning defined with respect to the command. Therefore, the intention of the command is not considered. Moreover, an erroneous input occurring due to a mistake in the operation by a user is not normally taken into account in the foregoing system.\nHitherto, a system is available in which a user is able to determine the level of the operational environment, such as whether help information, guidance information or the like from the system is provided, is determined by a user in accordance with his expertise with respect to the system or the complexity of the task to be performed in the system.\nThere has been another system in which a list of previously input commands is stored to which a reference can be made by a user or which enables a user to again execute the stored command.\nHowever, the foregoing conventional systems require a user to determine and change the level of the operational environment by operating buttons or inputting commands. Moreover, simple provision of the help information causes the help information having the same contents to be given to the user regardless of the expertise level of the user. Therefore, a skilled user is given unnecessarily detailed information, while the provided contents are too difficult for a beginner to be understood. Thus, the foregoing systems are unsatisfactory for the both types of users.\nThe conventional systems require a user to perform a complicated process to make a reference to the command list or again perform a command, thus requiring the user to be skilled in the procedure.\nAccordingly, an object of the present invention is to provide an information processing system capable of executing a process corresponding to input information and suitable to the situation and a method therefor.\nAnother object of the present invention is to provide an information processing system capable of controlling execution of a process in accordance with the attribute of a user and a method therefor.\nAnother object of the present invention is to provide an information processing system which makes a plan to meet a requirement from a user and which is capable of predicting a result of execution of the plan and a method therefor.\nAccording to one aspect, the present invention which achieves these objectives relates to an information processing system comprising: inputting means for inputting information; a knowledge base having knowledge of the relationship among situations, requirements corresponding to the situations and responses corresponding to the requirements; inference means which retrieves the knowledge base in accordance with the information input by the inputting means to infer a requirement corresponding to the information and the situation; and outputting means for outputting the requirement inferred by the inference means.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing system comprising: inputting means for inputting information; a knowledge base having knowledge of the relationship among a situation, a requirement, a response corresponding to the requirement, a plan corresponding to the requirement and a result of execution of a process; requirement inference means which retrieves the knowledge base in accordance with the information input by the inputting means to infer a requirement corresponding to the information; plan construction means for constructing a plan corresponding to the requirement; and predicting means for predicting a result of execution of the plan by performing a simulation of the plan made by the plan construction means.\nAccording to still another aspect, the present invention which achieves these objectives relates to an information processing system comprising: inputting means for inputting a command from a user; history storage means for storing input information performed by the inputting means; a knowledge base having knowledge of commands; inference means which infers the possible next-command to be input next by the inputting means the knowledge base in accordance with the information stored in the history storage means; and outputting means for outputting the possible commands inferred by the inference means.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing system comprising: inputting means for inputting information from a user; a knowledge base having knowledge of the attributes of each user, on the basis of the knowledge retrieved from requirements corresponding to the attributes of the users and responses corresponding to the requirements; inference means which retrieves the knowledge base in accordance with the information input by the inputting means and the attribute of the user, the information of which has been input, so as to infer a requirement which corresponds to the information and the attribute of the user; and output means for outputting the requirement inferred by the inference means.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing system comprising: natural language inputting means for inputting information in a natural language; a linguistic knowledge base having knowledge of languages; a global knowledge base having knowledge of common sense; analyzing means for analyzing, by using the linguistic knowledge base, the natural language information input by the natural language inputting means; a situation knowledge base having knowledge of the relationship among situations, requirements corresponding to the situations and responses corresponding to the requirements; requirement inference means which retrieves the situation knowledge base in accordance with the result of analysis performed by the analyzing means to infer a requirement corresponding to the natural language information; and outputting means for outputting the requirement inferred by the requirement inference means.\nAccording to another aspect, the present invention which achieves these objectives relates to an execution control system comprising: operation means which is operated by a user; executing means for executing a process in accordance with the operation performed by the operation means; a knowledge base having knowledge of operations to be performed by the operation means and the relationship among the operations and the attributes of the user; inference means for inferring the attribute of the user in accordance with the knowledge in the knowledge base; and control means for controlling execution to be performed by the executing means in accordance with the attributes of the user inferred by the inference means.\nAccording to another aspect, the present invention which achieves these objectives relates to an execution control system comprising: operation means to be operated by a user; executing means for executing a process in accordance with the operation performed by the operation means; history storage means for storing history of operations performed by the operation means; and control means for controlling execution of the process by the execution means in accordance with the history of the operations stored in the history storage means.\nAccording to another aspect, the present invention which achieves these objectives relates to an execution control system comprising: operation means to be operated by a user; executing means for executing a process in accordance with the operation performed by the operation means; history storage means for storing history of operations performed by the operation means; a knowledge base having knowledge of the operations of the operation means and the relationship between the operations and the attributes of the user; inference means for inferring the attribute of the use in accordance with the history of the operations stored in the history storage means and the knowledge in the knowledge base; and control means for controlling execution of the process by the executing means in accordance with the attributes of the user inferred by the inference means.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing method comprising the steps of: an inputting step for inputting information; an inference step which retrieves a knowledge base having knowledge of situations, requirements corresponding to the situations and the relationship between the requirements and responses in accordance with the information input in the inputting step to infer a requirement corresponding to the information and the situation; and an outputting step for outputting a requirement inferred in the inference step.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing method comprising the steps of: an inputting step for inputting information; a requirement inference step for retrieving, in accordance with the information input in the inputting step, a knowledge base having knowledge of the relationship among situations, requirements, responses corresponding to the requirements, plans corresponding to the requirements so as to infer a requirement corresponding to the information; a making means for making a plan corresponding to the requirement; and a predicting step for predicting a result of execution of the plan by performing a simulation of the plan made in the making step.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing method comprising the steps of: an inputting step for inputting a command from a user; a history storage step for storing information of history of inputs performed in the inputting step; an inference step in which a reference is made to a knowledge base having knowledge of commands in accordance with the information of history stored in the history storage step to infer possible commands to be next input in the inputting step; and an outputting step for outputting the possible commands inferred in the inference step.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing method comprising the steps of: an inputting step for inputting information from a user; an inference step for retrieving a knowledge base having knowledge of the relationship among the attribute of each user, requirements corresponding to the attributes of the users and responses corresponding to the requirements in accordance with the input information in the inputting step and the attribute of the user, who has input the information so as to infer a requirement which corresponds to the information and the attribute of the user; and an outputting step for outputting a requirement inferred in the inference step.\nAccording to another aspect, the present invention which achieves these objectives relates to an information processing method comprising the steps of: a natural language inputting step for inputting information in a natural language; an analyzing step for analyzing the natural language information input in the natural language inputting step by using a linguistic knowledge base having knowledge of languages and a global knowledge base having knowledge of common sense; a requirement inference step for retrieving the situation knowledge base having knowledge of the relationship among situations, requirements corresponding to the situations and responses corresponding to the requirements to infer a requirement corresponding to the natural language information; and an outputting step for outputting the requirement inferred in the requirement inference step.\nAccording to another aspect, the present invention which achieves these objectives relates to an execution control method comprising the steps of: an executing step for executing a process in accordance with the operation of an operation portion performed by a user; a knowledge base having knowledge of operations of the operation portion and knowledge of the relationship between the operations and the attributes of users; an inference step for inferring the attributes of the user in accordance with the knowledge of the knowledge base; and a control step for controlling execution of a process to be performed by the executing step in accordance with the attributes of the user inferred in the inference step.\nAccording to another aspect, the present invention which achieves these objectives relates to an execution control method comprising the steps of: an operation step in which an operation is performed by a user; an executing step for executing a process in accordance with the operation performed in the operating step; a history storage step for storing history of the operations performed in the operating step; and a control step for controlling execution of the process in accordance with the history of the operations stored in the history storage step.\nAccording to another aspect, the present invention which achieves these objectives relates to an execution control method comprising the steps of: an operation step in which an operation is performed by a user; an executing step for executing a process in accordance with the operation performed in the operating step; a history storage step for storing history of the operations performed in the operating step; a knowledge base having knowledge of the operations to be performed in the operation step and knowledge of the relationship between the operations and the attributes of the user; an inference step for inferring the attributes of the user in accordance with the history stored in the history storage step and the knowledge of the knowledge base; and a control step for controlling execution in the executing step in accordance with the attribute of the user inferred in the inference step.\nOther objectives and advantages besides those discussed above shall be apparent to those skilled in the art from the description of a preferred embodiment of the invention which follows. In the-description, reference is made to accompanying drawings, which form a part thereof, and which illustrate an example of the invention. Such example, however, is not exhaustive of the various embodiments of the invention, and therefore reference is made to the claims which follow the description for determining the scope of the invention."}
{"text": "The present disclosure relates to video coding.\nVideo coding finds use in many modern consumer electronic applications. Media players, such as media rendering applications, set top boxes and DVD players, commonly receive video data that has been coded by bandwidth compression techniques. The media players decode the video data before rendering it on a display. In other applications, videoconferencing applications perform bidirectional exchanges of coded video data. Each device codes video data representing locally-acquired video and transmits the coded video to another device. The other device receives and decodes the coded video, then renders it on a display.\nVideo coding and decoding processes typically are “lossy” processes. Video data recovered by decoders provides a representation of the source video from which it is derived but it possess various errors. When such errors are perceptible by human viewer, they often cause dissatisfaction with the viewing experience.\nIn many media exchange applications, different portions of a video may have very similar content. Consider, for example, a produced video where two characters are engaged in spoken dialogue with each other. Oftentimes, such events are represented by a video sequence that contains image information of a first speaker, then a second. Image content of the video sequence may toggle between image information of the two speakers for a time as the event progresses.\nIn many applications, video coders may apply different coding techniques at various points during such an event, which leads to different sets of artifacts. In the example above, a video coder may code image information of the first speaker differently during a first portion of spoken dialogue than during a second portion of dialogue, and the video coder may code representation of the first speaker differently during third, fourth, etc. portions as well. These different codings each may induce different sets of artifacts when the coded video data is decoded and rendered, which may lead to a dissatisfactory viewing experience."}
{"text": "1. Field of the Invention\nThe present invention relates to storage containers, and particularly, but not by way of limitation, to reusable containers for liquid additives for use in the oil well service industry.\n2. Description of the Prior Art\nMany activities conducted in the oil well servicing industry require the use of liquid materials which are typically provided in 55-gallon drums. There is a substantial cost involved in the handling and disposal of such drums which typically are not reused.\nThe prior art does include a number of liquid container systems designed for reuse.\nThe full-size liquid additive tank 26 disclosed herein is itself a part of the prior art in that it has been on sale by the assignee of the present invention for greater than one year prior to the filing of this application.\nA couple of other manufacturers have developed somewhat similar tanks such as those marketed by Hoover Universal, Inc., under the trademark AGRI-TOTE and such as the containers marketed under the trademark \"TransStore\" by Custom Metalcraft."}
{"text": "1. Field of the Invention\nThis invention relates generally to the directed self-assembly (DSA) of block copolymers (BCPs) using guiding lines to make an etch mask for pattern transfer into a substrate, and more particularly to a method for making the guiding lines.\n2. Description of the Related Art\nDirected self-assembly (DSA) of block copolymers (BCPs) has been proposed for making imprint templates for patterned media (also called bit-patterned media) magnetic recording disks and for patterning semiconductor devices, for example, for patterning parallel generally straight lines in MPU, DRAM and NAND flash devices. DSA of BCPs by use of a patterned sublayer that provides a chemical contrast pattern for the BCP film is well-known. After the BCP components self-assemble on the patterned sublayer, one of the components is selectively removed, leaving the other component with the desired pattern, which can then be used as an etch mask to transfer the pattern into an underlying substrate. The etched substrate can be used as an imprint template.\nIn conventional DSA by use of a chemical contrast pattern, an array of sparse guiding lines, usually made of a cross-linkable polymer mat, is lithographically patterned on the substrate. The pitch of these guiding lines (Ls) needs to be an integer multiple of the natural pitch (L0) formed by the block copolymer to be directed, i.e., Ls=nL0, where n is an integer equal to or greater than 1. However, the width of the guiding lines generally needs to be about 0.5 L0, or nearly equal to the width of one of the block copolymer components. While lithographic applications benefit from the fact that the guiding lines can be sparse (at a pitch of nL0) and that the high density comes from the density multiplication afforded by the block copolymers, the resolution requirement to pattern the width of the guiding lines keeps scaling with the final block copolymer dimensions. As lithographic applications advance to smaller dimensions, especially below 20 nm full pitch, the fabrication of guiding lines with a width below 10 nm while maintaining adequate roughness and width uniformity becomes ever more difficult and beyond what is possible with current resist materials for electron-beam (e-beam) or optical lithography.\nPending application Ser. No. 14/532,240 filed Nov. 4, 2014 and assigned to the same assignee as this application describes a pattern of guiding lines for DSA of a BCP that is an array of spaced guiding lines with oxidized sidewalls. A silicon-containing BCP self-assembles with the silicon-free BCP component wetting the oxidized sidewalls. The silicon-free BCP component is removed. The other BCP component then has a line density double that of the original spaced guiding lines and can be used as an etch mask to etch the underlying substrate. In the '240 application the spaced guiding lines are formed by electron-beam lithography. However, for small-pitch BCPs, it can be difficult to pattern guiding lines by e-beam lithography since it requires operating at the resolution limit of current e-beam resists. At the resolution limit of the e-beam resist, guiding lines are typically rough, broken, and/or wavy and cause defects in the BCP pattern that they guide. This effect is particularly severe for radial lines, which are required for making patterned media magnetic recording disks. Radial lines are more difficult to create by the e-beam tool due to the circumferential path of the beam, which results in an increased width roughness of the exposed lines. Also, radial lines require a smaller pitch than circumferential lines and thus a greater guiding line resolution to make patterned media bit structures with a high bit aspect ratio.\nWhat is needed is an improved method for making guiding lines that results in a BCP pattern after DSA with better quality and lower roughness."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device equipped with a light-emitting element or an electro-optical element in a pixel. In particular, the present invention relates to a display device having a layer including an organic material, a fluorescent material, or a phosphorescent material in the light-emitting element.\n2. Description of the Related Art\nA display device having a light-emitting element which includes a layer of an organic material between a pair of electrodes and emits light when a current is supplied between the electrodes has been developed. Such a display device has an advantage in reducing thickness and weight, has high visibility due to the self-luminance, and has high response speed. In addition, since power consumption of such a display device may potentially be made very small, it has been actively developed as a device of next generation, and some of such devices have been put into practical use.\nIn the display device using the light-emitting element having the aforementioned configuration, high image quality and widening of a color gamut have been expected. For example, a display device which can reproduce and display accurate colors in editing in a printing operation, seeing and listening to a work of art, movies, or the like, and catching a real color exactly in telemedicine, has been strongly expected. In view of this, in order to improve a color gamut which can be viewed by human eyes, research to optimize a structure such as improvement of color purity and widening of a color gamut has been carried out (for example, see Reference 1: Published Japanese translation of PCT International Publication for Patent Application No. 2001-039554).\nHowever, there is still a surplus in a color gamut which can be viewed by human eyes, and thus, a color reproduction area of a display device so far is still insufficient. FIG. 39 shows the CIE-XY chromaticity diagram which is established by COMMISSION INTERNATIONALE DE L'ECLAIRAGE (INTERNATIONAL COMMISSION ON ILLUMINATION: CIE) managing standards of color internationally. In an outer boundary of the diagram, a point which is near the rightmost end corresponds to an emission spectrum of 700 nm of red monochromatic light; a point which is near the uppermost end corresponds to an emission spectrum of 546.1 nm of green monochromatic light; and a point which is near the lowermost end corresponds to an emission spectrum of 435.8 nm of blue monochromatic light. In this chromaticity diagram, brightness (chroma) is lower in the inner side since the outer boundary of the graph (a visible area) corresponds to an emission spectrum of monochromatic light while the inner side thereof corresponds to a combination color obtained by combining different kinds of monochromatic light. In the case of expressing a color by an additive color mixture, a plurality of standard colors can reproduce only a color which lies in a position surrounded with a polygon formed of points which are shown in the CIE-XY chromaticity diagram.\nWhen red (R) is shown by the CIE-XY chromaticity diagram, human eyes can perceive a color which has a coordinate near a right region of the chromaticity diagram (a region surrounded with the circumference of the chromaticity diagram and a dotted line 3901 in FIG. 39) as red. In addition, when green (G) is shown by the CIE-XY chromaticity diagram, human eyes can perceive a color which has a coordinate near an upper region of the chromaticity diagram (a region surrounded with the circumference of the chromaticity diagram and a dotted line 3902 in FIG. 39) as green. When blue (B) is shown by the CIE-XY chromaticity diagram, human eyes can perceive a color which has a coordinate near a lower region of the chromaticity diagram (a region surrounded with the circumference of the chromaticity diagram and dotted lines 3903 and 3904 in FIG. 39) as blue. As a specific example, a hi-vision (high definition television broadcasting; HDTV) standard can be given, which has chromaticity coordinates of R (x=0.67, y=0.33), G (x=0.21, y=0.71), and B (x=0.14, y=0.08) (a triangle 3905 in FIG. 39).\nAccording to the method disclosed in Reference 1, a color reproduction area can be expanded in directions of arrows in FIG. 39 by increasing color purity, and brightness of colors recognized by human eyes can be increased. However, there is still a surplus in a color gamut which can be viewed by human eyes. Therefore, it is an essential task to expand the color reproduction area by satisfying the surplus in the color gamut which can be viewed by human eyes."}
{"text": "The deployment of an implant in the body is an often repeated and advantageous procedure performed during the practice of modern medicine. Implants may be used for innumerable purposes including the delivery of therapeutic for short and prolonged periods of time and for the structural support of failing internal cavities and lumens. Implants may be deployed in either case through non-invasive procedures such as endoscopy and through more invasive procedures that require large incisions into the body of a patient. The non-invasive and less-invasive procedures are generally used when the target area is accessible through a lumen of the body while the more invasive procedures may be employed when the target area is located deep within the body or otherwise not readily accessible through a lumen of the body.\nInvasive procedures have a higher financial cost and inflict a greater amount of trauma to the patient than non-invasive procedures. Consequently, non-invasive procedures are generally preferred over invasive procedures when given the choice. In certain circumstances, however, it may not be efficient or plausible to perform a procedure non-invasively. In these circumstances, more-invasive procedures may be required despite their shortcomings. For instance, when an implant needs to be deployed deep into the body in an area that is not readily accessible through a lumen or an external orifice an invasive procedure may be the only plausible technique. Similarly, when the target area is located deep within a muscle or tissue that does not have a nearby or sufficiently large lumen to provide access to the target site an invasive procedure may, again, be the only plausible technique. Consequently, in these situations as well as in others, a practitioner may be left with no choice but to perform an invasive procedure to deploy the implant to the targeted area."}
{"text": "1. Field of the Invention\nThis invention relates generally to a method for controlling liquid hydrocarbon spills.\nMore particularly, this invention relates to a method for bodying spilled hydrocarbon liquids to render those liquids non-flowing and amenable to recovery.\nIn one specific embodiment of the invention, oil spilled on water is contained and immobilized by dissolving in the oil particles of a relatively high molecular weight cryogenically comminuted, rubbery polymer applied to the oil as a slurry in a liquid cryogenic refrigerant.\n2. Description of the Prior Art\nA wide variety of solid materials have been used to absorb spilled liquid hydrocarbons. Materials commonly used include straw, newspaper, expanded perlite, and a number of polymeric materials of various compositions and configurations. Examples of polymers used to absorb and recover spilled hydrocarbons include particulate expanded polystyrene and polystyrenebutadione as shown by U.S. Pat. No. 3,929,631; a foam of an ethylene-alkyl acrylate copolymer recited in U.S. Pat. No. 3,819,514; granular polyurethane particles described in U.S. Pat. No. 3,657,125 and polyolefin fibers applied as a water slurry as shown by Japanese Patent No. 49-96980.\nIt is also known to apply refrigerants including solid or liquid carbon dioxide and liquid nitrogen to oil spills on water for the purpose either of solidifying the oil or of manipulating its movement. Cole et al in U.S. Pat. No. 3,614,873, for example, disclose the use of solid carbon dioxide particles of freeze oil floating on water thus allowing mechanical removal to the oil from the water. Ross et al in U.S. Pat. Nos. 4,031,707 and 4,129,431 disclose techniques for manipulating, or herding, a petroleum mass on water using blocks of solid carbon dioxide or other low temperature material to change the surface tension of the petroleum or to solidify it. Techniques utilizing refrigeration produce effects which necessarily are transitory in nature as both a water body and the atmosphere act as essentially infinite heat sinks which rapidly bring the chilled oil back to temperature equilibrium with the ambient environment.\nIt has also been recognized that a number of hydrocarbon polymers affect the viscosity characteristics of the hydrocarbon liquid in which they are dissolved. At low concentrations of polymer in hydrocarbon, in the general range of about 5 to 1000 ppm (0.0005% to 0.1%), the fluid flow friction loss is substantially reduced and this effect is utilized on a commercial basis to increasc flow of crude oil and petroleum products through pipelines. At somewhat higher polymer concentrations, on the order of 2000 to 5000 ppm (0.2% to 0.5%), a significant increase in viscosity of the hydrocarbon is observed. At and above about 5000 ppm the liquid begins to gel and at still higher polymer concentrations above about 2% depending upon the particular liquid hydrocarbon or hydrocarbon fraction, the liquid becomes a rubbery semi-solid or solid.\nIt is common knowledge in the art that non-crosslinked solid polymers are soluble to varying degrees in compatible liquids. The linear hydrocarbon polymers in general and the rubbery hydrocarbon polymers specifically are known to display a relatively high degree of solubility in hydrocarbon liquids. Examples of rubbery hydrocarbon polymers which display a significant degree of solubility in hydrocarbon liquids such as crude oil or kerosene include polyisoprene, polybutadiene, styrene-butadiene copolymers, polyisobutylene and the like.\nHowever, it is also recognized in the art that relatively high molecular weight polymers dissolve very slowly even in the best of solvents. Hours, days or even weeks of contact are required to dissolve such polymers in hydrocarbons at modestly high concentration levels. For example, Meier et al in U.S. Pat. No. 3,801,508 show that 22 hours of stirring contact was required to completely dissolve 1.5 grams of hydrogenated polyisoprene in 200 cc of cyclohexane. Dissolving the same polymer under identical conditions in West Texas crude oil required 70 hours. Attempts to dissolve ethylene-propylene copolymers in the same solvents were incomplete at the end of 168 hours.\nPrior art approaches to dissolving a polymer in a liquid hydrocarbon generally include contacting the polymer with the hydrocarbon at ambient to relatively elevated temperatures. It is generally accepted that an increase in temperature will speed dissolution of the polymer as will an increase in surface area as by comminution of the polymer into relatively small particles. Because many of the rubbery polymeric materials are relatively soft and resilient, they are extremely difficult and often impossible to comminute by conventional grinding techniques. Even after comminution, the soformed particles tend to stick and clump together thus negating the practical effect of comminution.\nOne example of prior art techniques for dissolving rubbery polymers; i.e., polyisobutylene, in hydrocarbons such as kerosene is shown by U.S. Pat. No. 3,215,154. Patentees disclose that a commercial polyisobutylene resin preground to a particle size of about 10 to 20 U.S. standard screen scale requires some 2 hours of stirring contact with kerosene until the solution attains sufficient viscosity to keep the undissolved resin particles dispersed. Another 3 to 5 hours is required for complete dissolution of the polymer particles. Patentees also teach that heating up to about 200.degree. F. increases the rate of solution.\nIt can thus be appreciated that, were prior art techniques used in the application of polymer particles to spilled hydrocarbons, little if any practical change in the physical characteristics of the spilled hydrocarbon could be expected."}
{"text": "1. Field of the Invention\nThe invention relates to the field of monitoring industrial process plants. More particularly, the invention relates to a method of process data mining.\n2. Description of the Prior Art\nIndustrial process plants employ various means to monitor processes and equipment. These may be Intelligent Electronic Devices (IED) specially made for particular equipment or some mathematical modeling method. IEDs are designed around the concept of a feedback control loop based on traditional control theory. Usually set points are defined within the range of data; violations of these set points trigger an alarm mode. These methods, while excellent at determining violation of set conditions and/or adjusting the control loop accordingly when system parameters change, are usually computationally intensive, even for monitoring and controlling a single piece of equipment. The complexity of such systems increases exponentially when implemented at the plant level, let alone for a fleet-wide implementation.\nOn the other hand, most mathematical methods employed in the industry currently utilize a two step process for analyzing and monitoring data: a training step and a monitoring step. In the training step, the system “learns” the modeling parameters, based on a-priori data, and in the monitoring step, the system monitors the equipment or the process against the “learned” result.\nThe major disadvantage of using such a two step process for learning and monitoring is that inconsistent results may be observed if the model changes due to a changing physical process, corrupt data, sensor malfunction etc. In such cases, there are no sure ways to know why the results are inconsistent and this gives rise to the possibility of false alarm.\nAnother obvious limitation is that a “clean” or optimal set of data is required to train the model before monitoring can be performed. This initial optimal data set usually is not available for the real world and has to be obtained either from simulation or from filtering existing prior data. If the model developed in such a way has some errors, or if some unforeseen event is seen in the data, the model has to be retrained to account for such a data point. This retraining of the model introduces an added effort in the monitoring step.\nWhat is needed, therefore, is a method for analyzing and classifying process data that does not require a clean set of data for training purposes, but that learns parameters on the fly by looking at changing data. What is further needed is such a method that optimizes the control data as more data is added. What is yet further needed is such a method that differentiates between key parameters and base parameters. What is still yet further needed is such a method that combines outputs of related sensors."}
{"text": "1. Field of the Invention\nThe present invention relates to a data recording disk, a disk drive device and a method for writing servo patterns on a data recording surface of a rotating data recording disk, and more particularly to a method, data recording disk and disk drive having servo patterns that improve efficiency of head positioning and data storage.\n2. Description of Related Art\nA giant magneto resistive (GMR) head has been developed as a next generation read/write head of a hard disk drive device. The GMR head can generate an output signal which is larger than an output signal generated by a prior magneto resistive (MR) head. The capability of generating the large output signal causes a reduction of a head width of the GMR head.\nThe GMR head is fabricated by the fabrication process of the MR head, and hence a head width of the GMR head is determined by a tolerance of the fabrication process of the MR head. It is assumed that the head width of the prior MR head is 2 xcexcm, (1 xcexcm=10xe2x88x926 m), the head width of the GMR head is 1 xcexcm, and the tolerance of the fabrication technology of the MR head is 0.4 xcexcm. Since the tolerance of the 0.4 xcexcm is applied to both the MR head and GMR head, the head width of the MR head is 2 xcexcmxc2x10.4 xcexcm, that is, the track width of the MR head varies from 1.6 xcexcm to 2.4 xcexcm, and a variance of the head width is 20%. The head width of the GMR head is 1 xcexcmxc2x10.4 xcexcm, that is, the head width of the GMR head varies from 0.6 xcexcm to 1.4 xcexcm, and a variance of the head width is 40%. In this manner, the variance of the head width of the MR head is 20%, while the variance of the head width of the GMR head is 40% which is larger than that of the MR head.\nAn uniform track pitch is used in the track arrangement of the prior MR technology. A track pitch is a distance between a track center of one track and a track center of the next track.\nIf the uniform track pitch of the MR head is used for the GMR head, the following first problem, as shown in FIG. 1 occurs. As described above, the head width of the GMR head varies from 0.6 xcexcm to 1.4 xcexcm. For simplifying the description, three kinds of the GMR heads 1, 2 and 3 are shown in FIG. 1, i.e., the GMR head 1 has the 1.4 xcexcm head width, the GMR head 2 has the 1.0 xcexcm head width, and the GMR head 3 has the 0.6 xcexcm head width. The track width of all the tracks 1, 2 and 3 are the same width, i.e., 1.4 xcexcm, which is the largest width of the GMR head, and the uniform track pitch is used for all the tracks. A gap is provided for separating the adjacent two tracks. And, the center of the GMR heads 1, 2 and 3 are aligned to the track center of each track, respectively. It is apparent that useless or wasted space exists between the GMR heads 1 and 2, and between the GMR heads 2 and 3.\nThe following second problem occurs with respect to the servo patterns provided in the track. The servo patterns A and B are written to define the tracks, and read by the head to generate a feedback signal which is supplied to a hard disk control device. The detail of the servo patterns A and B are shown in FIG. 4(B) and (C). The servo pattern A has alternately arranged opposite magnetization direction, and the servo pattern B has an uniform magnetization direction. The hard disk control device responds to the feedback signal to position the center of the head on the track center. It is assumed that the GMR heads 1, 2 and 3 are shifted from the track center by a distance of 10% of the track pitch, as shown by the dashed line in FIG. 1. For example, it is assumed that the track pitch=track widthxc3x971.1=1.4 xcexcm=1.1xc3x971.6 xcexcm, and the shift distance=track pitchxc3x970.1=0.16 xcexcm. In this case, the GMR heads 1, 2 and 3 are shifted with respect to the servo pattern B, and the variation of the levels of the feedback signals sensed by the GMR heads 1, 2 and 3 are as follows.\nIn the case of the GMR head 1 of the 1.4 xcexcm head width: 0.16 xcexcm/1.41 xcexcm=0.114. That is, the level of feedback signal is reduced by 11.4%. In the case of the GMR head 2 of the 1.01 xcexcm head width: 0.16 xcexcm/1.0 xcexcm=0.16. That is, the level of feedback signal is reduced by 16%. In the case of the GMR head 3 of the 0.6 xcexcm head width: 0.16 xcexcm/0.6 xcexcm=0.267. That is, the level of feedback signal is reduced by 26.7%.\nThe variation of the levels of the feedback signal sensed by the prior MR head of 2.4 xcexcm head width, the prior MR head of 2.0 xcexcm head width, and the prior MR head of 1.6 xcexcm head width are, as follows. It is assumed that the track pitch=track widthxc3x971.1=2.4 xcexcmxc3x971.1=2.6 xcexcm, and the shift distance=track pitchxc3x970.1=0.26 xcexcm.\nIn the case of the MR head of the 2.4 xcexcm head width: 0.26 xcexc m/2.4 xcexc m=0.108. That is, the level of the feedback signal is reduced by 10.8%. In the case of the MR head of the 2.0 xcexcm head width: 0.26 xcexcm /2.0 xcexcm=0.13. That is, the level of the feedback signal is reduced by 13%. In the case of the MR head of the 1.6 xcexcm head width: 0.26 xcexcm/1.6 xcexcm=0.163. That is, the level of the feedback signal is reduced by 16.3%.\nIn this manner, when the head is shifted by the distance of track pitch multiplied by A, the level of the feedback signal of the MR head to the hard disk control device varies from 10.8% to 16.3%, while the level of the feedback signal of the GMR head to the hard disk control device varies from 11.4% to 26.7%.\nThe range of the variation of the feedback signals to the hard disk control device should be small, since if the range of the variation becomes large, undesirable phenomenon, such as an oscillation of the entire servo system may occur.\nIt can be seen then that there is a need for a method, data recording disk and disk drive that eliminates useless or wasted space on a data recording disk.\nIt can also be seen that there is a need for a method, data recording disk and disk drive that provides a track pitch that varies depending o the track width.\nIt can also be seen that there is a need for a method, data recording disk and disk drive that provides a uniform level of variation I the feedback signals.\nTo overcome the limitations in the prior art described above, and to overcome other limitations that will become apparent upon reading and understanding the present specification, the present invention discloses a data recording disk that includes a first data recording surface and a second data recording surface, wherein a track pitch of adjacent data recording tracks of the first data recording surface differs from a track pitch of adjacent data recording tracks of the second data recording surface.\nA data recording disk according to the present invention includes a first data recording surface and a second data recording surface, wherein the number of data recording tracks of the first data recording surface differs from the number of data recording tracks of the second data recording surface, and the a gap width between the data recording tracks of the first data recording surface differs from a gap width between the data recording tracks of the second data recording surface.\nA data recording disk according to the present invention includes a first data recording surface and a second data recording surface, wherein the number of data recording tracks of the first data recording surface differs from the number of data recording tracks of the second data recording surface.\nA disk drive device according to the present invention includes data recording surfaces on a rotating data recording disk, each of the data recording surfaces is provided with a read/write head, wherein the number of data recording tracks of a first data recording surface differs from the number of data recording tracks of a second data recording surface, and a gap width between the data recording tracks of the first data recording surface differs from a gap width between the data recording tracks of the second data recording surface.\nAll data recording tracks of the first data recording surface have substantially the same track width, and all gaps of the first data recording surface have substantially the same gap width.\nAll data recording tracks of the second data recording surface have substantially the same track width, and all gaps of the second data recording surface have substantially the same gap width.\nData is written in a direction from an outer most data recording track to an inner most data recording track of the first data recording surface, then the data is written in a direction from an outer most data recording track to an inner most data recording track of the second data recording surface.\nAddresses of data recording tracks are successively assigned from an outer most data recording track to an inner most data recording track of the first data recording surface, and the address of the inner most data recording track of the first data recording surface is followed by an address of an outer most data recording track of the second data recording surface, and the addresses on the second data recording surface are successively assigned from the outer most data recording track to an inner most data recording track of the second data recording surface.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on a rotating data recording disk by a read/write head contained in a disk drive device, includes positioning the read/write head at one position along a radial direction of the rotating data recording disk, writing servo patterns on the rotating data recording disk by the read/write head positioned at the one position, reading the servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50xe2x88x92Gw/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Gw is a predetermined gap width between the data recording tracks, stopping the read/write head when the detection is made, and overlappingly writing the servo patterns on the rotating data recording disk by the stopped read/write head.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on a rotating data recording disk by a read/write head contained in a disk drive device, includes specifying a gap width between adjacent data recording tracks, writing servo patterns for defining a plurality of data recording tracks separated by the specified gap width on the rotating data recording disk by the read/write head, detecting that the total number of data recording tracks is smaller than a predetermined number, specifying an overlap width of adjacent data recording tracks, and rewriting servo patterns for defining a plurality of data recording tracks overlapped each other by the overlap width on the rotating data recording disk by the read/write head, the rewriting of servo patterns including positioning the read/write head at one position along a radial direction of the rotating data recording disk, rewriting the servo patterns on the rotating data recording disk by the read/write head positioned at the one position, reading the rewritten servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the rewritten servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50+Y/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Y is the overlap width, stopping the movement of the read/write head when the detection is made, and rewriting the servo patterns on the rotating data recording disk by the stopped read/write head stopped.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on each of a plurality of data recording surfaces on rotating data recording disks by a read/write head provided on each of the plurality of data recording surfaces contained is a disk drive device, includes specifying a gap width between adjacent data recording tracks, writing servo patterns on a plurality of servo tracks for defining the plurality of data recording tracks separated by the specified gap width on one data recording surface by its read/write head, storing a count value indicating the number of servo tracks written on the one data recording surface, repeating the writing of servo patterns and storing the count value until the servo patterns are written on all data recording surfaces, determining whether the total number of data recording tracks written on the data recording surfaces exceeds a predetermined number, or not, specifying an overlap width of adjacent data recording tracks, if the total number of data recording tracks written on data surfaces does not exceed the over the top width, finding the data recording surface containing the least number of data recording tracks by determining the count values of the all data recording surfaces, and rewriting the servo patterns for defining the plurality of data recording tracks overlapped each other by the overlap width on the data recording surface found to contain the least number of data recording tracks to its read/write head, the rewriting including positioning the read/write head at one position along a radial direction of the data recording surface, rewriting servo patterns on the data recording surface by the read/write head positioned at the one position, reading the rewritten servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the rewritten servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50+Y/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Y is the overlap width, stopping the movement of the read/write head when the detection is made, and rewriting the servo patterns on the data recording surface by the stopped read/write head stopped.\nA method according to the present invention for writing servo patterns for defining a plurality of data recording tracks on each of a plurality of data recording surfaces on rotating data recording disks by a read/write head provided on each of the plurality of data recording surfaces contained is a disk drive device, includes specifying a gap width between adjacent data recording tracks, positioning the read/write head at one position along a radial direction of one data recording surface, setting a count value of a counter to an initial value, writing servo patterns on the data recording surface by the read/write head positioned at the one position, reading the servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50xe2x88x92Gw/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Gw is the specified gap width between the data recording tracks, stopping the read/write head when the detection is made, incrementing the count value, overlappingly writing the servo patterns on the data recording surface by the stopped read/write head, repeating the reading, moving, reading and stopping until the servo patterns are entirely written on the one data recording surface, electing next data recording surface and repeating the positioning, setting, writing, reading, moving, reading, stopping, incrementing, overlapping and repeating, determining whether the total number of data recording tracks written on all data recording surfaces exceeds a predetermined number, or not, specifying an overlap width of adjacent data recording tracks, if the total number of data recording tracks written on all data recording surfaces does not exceed the predetermined number, finding the data recording surface containing the least number of data recording tracks by determining the count values of the all data recording surfaces, and rewriting the servo patterns for defining the plurality of data recording tracks overlapped each other by the overlap width on the data recording surface found containing the least number of data recording tracks by its read/write head, the rewriting including positioning the read/write head at one position along a radial direction of the data recording surface, rewriting servo patterns on the data recording surface by the read/write head positioned at the one position, reading the rewritten servo patterns by the read/write head positioned at the one position to detect an amplitude of output signal generated by the read/write head, moving the read/write head from the one position along the radial direction, reading the rewritten servo patterns by the read/write head moved from the one position to detect that an amplitude of output signal generated by the moved read/write head becomes 50+Y/2% of the amplitude of output signal generated by the read/write head when it is positioned at the one position, wherein the Y is the overlap width, stopping the movement of the read/write head when the detection is made, and rewriting the servo patterns on the data recording surface by the stopped read/write head stopped.\nThe total number of data recording tracks written on the all data recording surfaces is determined by referring to the count values of the all data recording surfaces.\nThese and various other advantages and features of novelty which characterize the invention are pointed out with particularity in the claims annexed hereto and form a part hereof. However, for a better understanding of the invention, its advantages, and the objects obtained by its use, reference should be made to the drawings which form a further part hereof, and to accompanying descriptive matter, in which there are illustrated and described specific examples of an apparatus in accordance with the invention."}
{"text": "This invention relates to a slider block for a scroll compressor wherein the inner bore of the slider block has a clearance recess to allow for movement of the eccentric pin that moves relative to the slider block bore.\nScroll compressors are becoming widely utilized in refrigerant compression applications. In a scroll compressor a first scroll member has a generally spiral wrap extending from a base. A second scroll member has its own wrap which interfits with the spiral wrap of the first scroll member. The second scroll member is caused to orbit relative to the first scroll member to entrap and then compress a refrigerant. The second scroll member is generally driven to orbit by an electric rotary motor driving the second scroll member through a Oldham coupling. The connection between the driveshaft of the motor and the orbiting scroll is through a slider block, such that the second scroll has the ability to move relative to the first scroll under certain circumstances. Typically, the shaft has an eccentric pin extending upwardly into an opening in the slider block. The opening has a flat drive surface that is in contact with a barrel shaped drive surface on the eccentric pin. The bore of the slider block has two opposed curved surfaces. The pin may slide relative to the flat surface on the slider block, such that the orbiting scroll can move towards and away from the wrap on the first scroll member.\nIn one type of slider block, the accurate surfaces that define the two curved portions of the bore need to be spaced relatively far inwardly to provide for additional structure on the outer periphery of the slider block. In particular, one type of slider block has an oil notch at its side. The formation of the oil notch requires that the curved surface be spaced relatively far inwardly to provide for sufficient wall thickness. However, if the wall is sufficiently thick such that the curved surface is spaced inwardly, it may well be there is insufficient clearance for the pin to move to certain positions with out contacting the opposed curved surface.\nThe present invention is directed to addressing the above-mentioned problem.\nA scroll compressor is provided with a slider block having a pair of curved opposed surfaces. The surfaces are preferably centered on a common center line. A pair of drive flats separates the two curved surfaces. A recess is formed into one of the two curved surfaces to provide additional clearance for movement of the eccentric pin. Preferably the recess is generally opposed to structure on the outer periphery of the slider block; the structure is preferably a notch to allow for flow of oil. The notch is associated with one of the two curved surfaces, and the curved surface is spaced inwardly from the notch to provide sufficient wall thickness. The inwardly spaced curved surface results in the need for additional clearance provided by the recess.\nThese and other features of the present invention can be best understood from the following specification and drawings, the following of which is a brief description."}
{"text": "The demand for microphones for mobile equipment such as mobile telephones, headsets and cameras tends to follow the growing demand for mobile equipment—for instance mobile telephones.\nThe demand has for many years been rather simple in that the demand was for microphones with extremely low costs and for microphones suitable for production in very high volumes. The performance of such microphones was comparable from one manufacturer to another and was at a level comparable to that of telephony systems. In recent years, however, the demand has changed to also be for microphones with a performance above that of telephony systems. Today there appears to be a trend in the demand heading towards so-called high-fidelity (hi-fi) quality.\nUse of integrated digital processors, with ever increasing performance, in the different types of mobile equipment has also brought attention to the performance of the more peripheral links of the signal processing chain from pick-up of a signal over transmission and/or storage to reproduction of the signal. Such more peripheral link is for instance the microphone or the circuitry embodied with a microphone transducer in a microphone capsule. A microphone capsule—also denoted a microphone element—can include shock mounts, acoustic isolators, protective covers and a semiconductor die, with integrated circuitry, in addition to the microphone transducer. The microphone transducer and the integrated electronic circuit is embodied in the microphone capsule converts acoustic energy to electrical energy so as to provide an electrical microphone signal.\nIt has been discovered that integrated digital processors can be configured to repair some damages to a microphone signal occurred due to inadequate signal conditioning in the microphone capsule. But in general it is far more efficient not to disregard aspects of signal conditioning in peripheral links of the signal processing chain to thereby avoid destroying the microphone signal and consequently being able to provide far better repair of the signal if needed at all. The microphone signal may be destroyed by disregarding a noise source and/or by overloading an amplifier (in the capsule).\nThus, there is a demand for high quality microphones, but unfortunately the demand for low price seems to persist. As cost of a semiconductor die is directly related to the size of the die it is important, for the purpose of reducing price, that the electronic circuit integrated on the die is as small or compact as possible. Therefore, very simple circuits are desired, with a due regard to a desired (high) performance.\nIt has been discovered that meeting high performance demands is not simply a question of providing a more robust or conservative design. Due to the important cost issues and signal conditioning aspects it is found that there is no single one fixed signal conditioning circuit that is able to provide high performance in various acoustic situations. Such various situations could be described as a voice signal with/without loud/quiet background noise, a loud or quiet voice signal or combinations thereof. Thus, the signal conditioning needed to provide high performance is different from one situation to another.\nIt has therefore been proposed—despite the additional cost of a more complex semiconductor die in the capsule—to provide the semiconductor die with means for the circuit on the die to adapt to a given acoustic situation. Thereby, high performance can be achieved in different acoustic situations. In some designs of microphones it may be equitable to provide the adaptation to different acoustic situations by a control loop embodied entirely on the semiconductor die, but in other designs it may be equitable to provide the control loop to provide a control feedback from a circuit external to the microphone capsule. Thereby it is necessary to configure the semiconductor die for an external circuit to make the circuit in the semiconductor die adapt to a situation. To this end cost is generally an obstacle simply to have one or more additional pad(s) for receiving such an external feedback.\nThus, since high quality microphones are sought after, more complex circuitry is inherently needed which—all other things being equal—has a higher power consumption. Since the mobile equipment is battery powered current consumption of the device including portion thereof is subject to be minimized as far as possible. This adds an additional and important dimension to the demand.\nThe microphone is based on the principle of a capacitor, which is formed by a movable member that constitutes a membrane of the microphone and another member, e.g. a so-called back plate of the microphone. One of the members of the microphone, preferably the membrane, is charged by a constant electrical charge. The charge is either provided as an electrostatic charge captured on one of the members or provided by a voltage source e.g. a charge pump or voltage step-up circuit on the semiconductor die.\nA sound pressure detected by the microphone will cause the membrane to move and consequently change the capacitance of the capacitor formed by the membrane member and the other member. When the charge on the capacitor formed by these two members is kept constant, the voltage across the two capacitor members will change with the incoming sound pressure level. As the charge on the microphone capacitor has to be kept constant to maintain proportionality between sound pressure and voltage across the capacitor members, it is important not to load the microphone capacitance with any resistive load. A resistive load will discharge the capacitor and thereby degrade or ruin the capacitors performance as a microphone. A capacitive load will reduce sensitivity of the microphone transducer.\nTherefore, in order to pick up a microphone signal from the capacitor, amplifiers configured with the primary objective of providing high input resistance are preferred to buffer the capacitor from circuits which are optimized for other objectives. The amplifier connected to pick up the microphone signal is typically denoted a preamplifier or a buffer amplifier or simply a buffer. The preamplifier is typically connected physically very close to the capacitor—within a distance of very few millimetres or fractions of millimetres.\nFor small sized microphones only a very limited amount of electrical charge can be stored on one of the microphone members. This further emphasizes the requirement of high input resistance. Consequently, the input resistance of preamplifiers for small sized microphones has to be extremely high—in the magnitude of Giga ohms. Additionally, the input capacitance of this amplifier has to be very small in order to achieve a fair sensitivity to sound pressure.\nTraditionally, this preamplifier has been implemented as a simple JFET. The JFET solution has been sufficient, but demands in the telecom industry call for ever smaller microphones—with increased sensitivity. This yields a contradiction in terms since sensitivity of the microphone capacitor drops as size goes down. All other things being equal, this will further reduce the sensitivity of the microphone and the buffer in combination. The demands in the telecom industry are among other things driven by market trends which encompass hands free operation of different types of small-sized equipment and more widespread application of microphones in e.g. camera applications.\nSo obviously, there is a need for microphone preamplifiers with gain and very low input capacitance, and lowest possible preamplifier die area. Additionally, low noise is important. Low noise is important as noise, during design of the microphone, can be traded for area—i.e. if the circuit has low noise and a noise lower than required, this noise level overhead can be traded for lower chip die area and it is thus possible to manufacture the preamplifier at lower cost.\nWhen designing a preamplifier for a microphone there is normally three noise sources. These sources are noise from a bias resistor, 1/f noise from an input transistor, and white noise from the input transistor. Typically, input transistor noise dominates. Both white noise and 1/f noise can be minimized by optimizing the length and the width of the input transistor(s). This applies for any input stage e.g. a single transistor stage or a differential stage.\nThe noise from the bias resistor can also be minimized. If the bias resistor is made very large then the noise from the resistor will be high pass filtered and the in-band noise will be very low. This has the effect though that the lower bandwidth limit of the amplifier will be very low. This can be a problem as the input of the amplifier will settle at a nominal value only after a very long period of time after power up. Additionally, signals with intensive low frequency content arising form e.g. slamming of a door or infra sound in a car can overload the amplifier. Another related problem is small leakage currents originating from mounting of the die inside a microphone module. Such currents will, due to the extreme input impedance, establish a DC offset. This will reduce the overload margin of the amplifier.\nMicrophones is typically manufactured with a yield of approximately 80-90% i.e. 80-90% of the total number of produced microphones satisfy specifications on their performance. Unfortunately, 10-20% of the production is discarded since for instance sensitivity of the microphone does not satisfy the specification. A solution to reducing the discarding rate would be highly appreciated by the industry.\nAnother problem of e.g. electret microphones is the ageing phenomena's in which the electret microphone might change sensitivity over time thereby leading to a discrepancy between the electret microphone output and the gain of the buffer amplifier.\nA microphone subjected to a background noise comprising low frequency sound at high amplitudes, e.g. from a motorized vehicle, may be prone to the problem of e.g. clipping the sound signal from the microphone. In case a voice signal is present in combination with such a background signal, the information in the voice signal may be lost since the sound pressure results in the corresponding electrical signal being clipped. The clipping of the microphone signal may occur when the amplitude of the low frequency background superposed on the voice signal overload the amplifier amplifying the signal from the microphone e.g. by exceeding a maximum sound pressure that the microphone and amplifier may handle e.g. 110 dB SPL. Minor overloading of the amplifier may result in signal clipping while severe overloading of the amplifier may yield a period of time, e.g. in the order of seconds, where the amplifier has ceased to operate as an amplifier."}
{"text": "The invention relates to a method for setting a conicity of a die of a strand casting installation, in particular a slab strand casting installation, during a casting procedure, comprising the following method steps: measuring temperature values along at least one centric measuring path running in a casting direction along an adjustably disposed die wall; and measuring temperature values along at least one peripheral measuring path running in a casting direction along the adjustably disposed die wall, wherein the peripheral measuring path runs between the centric measuring path and a lateral periphery of the die wall, and a spacing of the peripheral measuring path from said lateral periphery of the die wall is smaller than a spacing of the centric measuring path from the other lateral periphery of the die wall.\nThe invention furthermore relates to a device for a strand casting installation, in particular a slab strand casting installation, comprising: at least one die having mutually opposite, adjustably disposed die walls; at least one adjustment mechanism for adjusting the die walls; at least one actuation electronics system for actuating the adjustment mechanism; at least one centric sensor installation, connected to the actuation electronics system, for measuring temperature values along at least one centric measuring path running in a casting direction along one of the adjustably disposed die walls; and at least one peripheral sensor installation, connected to the actuation electronics system, for measuring temperature values along at least one peripheral measuring path running in a casting direction along the adjustably disposed die wall, wherein the peripheral measuring path runs between the centric measuring path and a lateral periphery of the die wall, and a spacing of the peripheral measuring path from said lateral periphery of the die wall is smaller than a spacing of the centric measuring path from the other lateral periphery of the die wall.\nA respective method and a respective device are known from DE 10 2014 227 013 A1.\nIn the strand casting of metals in a strand casting installation, the strand width of a cast strand can be set by way of adjustable die walls on mutually opposite narrow sides of a die. The die comprises a casting passage that continues from an inlet end by way of which a metal melt is cast into the die, to an outlet end from which a strand having a strand shell and a liquid core exits, said casting passage in the cross section perpendicular to a casting direction which is defined as the direction of the metal melt flowing through the die having a rectangular shape. The strand can be configured as a slab, for example.\nDuring cooling of the metal melt within the die, a loss of volume arises, this meaning that a cross-sectional area of the metal melt, or of the strand created therefrom, respectively, is larger at the inlet end of the die than a cross-sectional area of the metal melt, or of the strand, respectively, at the outlet end of the die. In order for a desired guiding of the metal melt, or of the strand, respectively, to be guaranteed by way of a contact with the die across the full area, the casting passage of the die is configured so as to conically taper in the casting direction. The conicity of the casting passage of the die is usually set by way of a variation in the incline of the die walls adjustably disposed on mutually opposite narrow sides. Without such a conical configuration of the casting passage of the die, the strand would no longer be guided in the lower region of the die, or in the region of the outlet end, respectively. In this instance, a controlled dissipation of heat from the strand to the usually water-cooled die walls is no longer possible. For the sake of simplicity, the conicity of the casting passage of the die in the context of the present application is referred to as the conicity of the die.\nFIG. 1A shows a result in the case of an optimal conicity of a die 1. Shown is a lateral portion of the die 1 in the cross section perpendicular to the casting direction, said die 1 having two mutually opposite die walls 2 which are in each case assigned to one broad side of the die 1. Moreover, the die 1 has two mutually opposite, adjustably disposed eyeballs 3 which in each case are assigned to a narrow side of the die 1, and of which only one die wall 3 is shown in FIG. 1A. A strand 4 which has a strand shell 5 and a core 6 of liquid metal is guided through the die 1. The strand 4 bears on the entire area of the die walls 2 and 3, and on account thereof is guided in an optimal manner and simultaneously cooled by means of the die walls 2 and 3.\nFIG. 1B shows a result in the case of an inadequate conicity of a die 1. On account thereof, the strand 4, or the strand shell 5, respectively, bulges outwards on the narrow sides, and the strand 4 simultaneously flattens towards the periphery, thus having a lesser thickness at said periphery than in a central region. Moreover, the periphery of the strand 4 no longer bears on the die walls 2 and 3 across the entire area such that the strand 4 is not guided in an optimal manner and is non-uniformly cooled by means of the die walls 2 and 3.\nFIG. 1C shows a result in the case of an excessive conicity of a die 1. On account thereof, a concavity 6 arises on a narrow lateral face of the strand 4, or of the strand shell 5, respectively, and longitudinal depressions 7 near the periphery, the so-called rain gutters, arise on the broad sides of the strand 4. On account thereof, the narrow sides and the peripheral portions of the broad sides of the strand 4 no longer bear on the die walls 2 and 3 across the entire area such that the strand 4 is not guided in an optimal manner and is non-uniformly cooled by means of the die walls 2 and 3. On account of the inadequate cooling on the narrow sides of the strand 4, the strand shell 5 on the narrow sides of the strand 4 is not configured so as to be sufficiently strong. On account thereof, bulges are formed on the narrow sides of the strand 4 after the strand 4 has exited the die 1 and is in a secondary cooling in a strand casting installation.\nThe formation of concavities on the narrow sides of a strand can conventionally be minimized or eliminated, respectively, by reducing the conicity of the die. By contrast thereto, if it is established that the strand after passing through the die has bulges on the narrow sides, the conicity of the die is usually increased. However, if the bulges on the narrow sides of the strand are caused by an excessive conicity of the die according to FIG. 1C, the defect pattern of the longitudinal depressions on the broad sides of the strand are amplified on account of the further amplification of the conicity of the die for the intended reduction of the bulges, as is shown in FIG. 1D. This can lead to fissures on the strand shell on the base of the respective longitudinal depression, and under certain circumstances can lead to strand breakages in the region of the respective longitudinal depressions. The overall quality of the strand is reduced on account of the longitudinal depressions."}
{"text": "This invention pertains to a circuit for multiplying one variable by another and is particularly concerned with a circuit which is useful in a number of applications one of which is in an electronic engine control circuit for use with an internal combustion engine spark timing system.\nIn U.S. Pat. No. 3,978,833, dated Sept. 7, 1976, and assigned to the same assignee as the present application, there is disclosed a circuit for providing a programmed engine control function signal which is specifically disclosed as a programmed vacuum advance signal. This programmed vacuum advance signal is developed from engine manifold vacuum and from an accumulation function which is representative of the duration for which the engine has been operating in selected operating modes. Within its control range between maximum and minimum limits of vacuum spark advance, the programmed vacuum advance signal is developed as a function of the product of these two variables, namely (1) instantaneous engine manifold vacuum, as monitored by a vacuum transducer and associated transducer circuit, and (2) the accumulation function which increases in magnitude as a function of time when the engine is in a non-idle operating condition and which decreases in magnitude as a function of time when the engine is operating at idle condition. The vacuum transducer and associated transducer circuit disclosed in that patent develop a constant frequency rectangular waveform wherein the pulse width of the pulses comprising the waveform is representative of engine manifold vacuum. However, the relationship is such that at zero manifold vacuum the pulses have a definite non-zero pulse width. Thus, the product of the pulse width of the rectangular waveform signal by the accumulation function would yield an incorrect result where it is desired to obtain the product of manifold vacuum by the accumulation function. In order to obtain the correct product of manifold vacuum and accumulation function signal, it becomes necessary to make allowance for the non-zero width of the pulses which exists at zero manifold vacuum. Accordingly, a special circuit which makes this allowance is provided in that patent.\nThe present invention is concerned with an improvement upon the scheme disclosed in the foregoing patent. The advantages of the improvement afforded by the present invention are as follows.\nThe gain of the output amplifier is much easier to set. The previous system uses a more complex vacuum amplifier requireing proper choice of four interactive resistors. It is often difficult to find four appropriate resistors among standard values in order to achieve a desired gain while staying within impedance limits imposed by the system. With the present invention, development of engine control systems can be speeded up as systems can now be provided with a variable gain adjusted by a potentiometer. This capability was not available in the previous system as simultaneous adjustment of four interacting resistors was not practical. Additional advantages of the invention will be seen as the description proceeds."}
{"text": "This invention relates to a speaker supporting unit which is most adapted to attach a speaker unit to a door panel of an automobile.\nIn an audio system installed in an automobile, a speaker unit is often attached to a door panel of the automobile. One of the prior art technique for mounting the speaker unit is shown in FIG. 12 and FIG. 13 which employs a pair of dish-shaped mounting bases 62 and 64. The interior mounting base 62 is inserted through a circular opening 60a formed in the door panel 60 to be engaged within the other mounting base 64. These mounting bases 62 and 64 are then joined together by fastening screws 65, 65. The mounting base assembly is thus secured to the door panel 60 by rims thereof tightened across the door panel 60. There are interposed between the mounting bases 62 and 64 a predetermined number of washers 66 corresponding to the thickness of the door panel 60, to prevent deformation of the mounting bases when tightened. A speaker 70 to which a grille lid 68 is attached is contained in the mounting base 62 and secured by any suitable fastening means (not shown). In this prior art, a spacer 72 may be additionally employed in a manner as particularly shown in FIG. 13 to change and adjust the orientation of the speaker 70, in which case the screw 74 should be accompanied by a bracket 76 to thereby extend in a direction perpendicular to the inclined orientation of the speaker 70.\nAccording to the foregoing prior art, however, the orientation of the speaker 70 will be determined in dependence upon the inclination of the spacer 72. It has been impossible to adjust as desired the orientation of the speaker within a wide adjustable range. One or more sets of the spacer and the corresponding bracket should be prepared and selectively used upon necessity, which not only is troublesome but also may cause accidental loss of the adjustment members. Further, as the inner diameter of the mount base should be large enough to allow the orientation of the speaker to be changed, there might be a great probability of foreign matter entering into the mount base."}
{"text": "The present disclosure relates to an optical communication device, a reception apparatus, a transmission apparatus, and a transmission and reception system.\nWith the recent development of information-oriented society, an information processing apparatus such as personal computers (PCs) and servers handles an explosively increasing amount of information (the amount of data or signal). With such an increase in the amount of data, there is a growing necessity to transmit and receive more data between apparatus, devices, and chips at higher speed.\nAs an example, in a processor such as central processing units (CPUs) used in computers, the increase in its speed and functionality is progressing. In related art, the transmission and reception of data between processors are mainly carried out, for example, by using electrical signals through copper wiring on a substrate. However, the high-speed data transmission carried out by using such data transmission technology employing electrical signals (electrical communication or interconnect technology) causes problems of RC signal delay, impedance mismatching, EMC/EMI, and crosstalk, which makes it difficult to further increase data transmission speed.\nThus, the optical communication technology that modulates electrical signals into optical data and transmits data using light (or optical interconnect technology) has been developed as data transmission technology that can be substituted for the electrical communication technology. For example, JP 2005-181610A discloses a socket for optical communication that includes an installation part for installing an optical waveguide, a light emitting element for making light incident to the optical waveguide, and a light receiving element for receiving light emitted from the optical waveguide. In this socket, at least one of the light emitting element and the light receiving element is arranged corresponding to the optical waveguide. In addition, JP 2007-25310A discloses a technique that makes a socket for optical communication thinner by defining a recess in the socket and fitting an interposer substrate into the recess."}
{"text": "A previous invention disclosed in U.S. patent Ser. No. 13/089,576 describes the production of a concentrated organic soil amendment. This organic soil amendment and similar concentrated compost tea is typically supplied in a concentrated liquid form. This concentrated organic soil amendment allows reduced shipping volume and weight of the product. It also can be stored at reduced temperature thereby minimizing adverse effects to the efficacy of the product. The concentrated liquid soil amendment or fertilizer needs to be diluted before administering to plants and/or soil, such as turf grass or crops. When the administering area is large, a vehicle mounted tank is typically required. Diluting large quantities of concentrate can be slow and cumbersome, and may also pose occupational hazard to the operators if they have to carry heavy containers of the soil amendment or fertilizer atop a vehicle mounted tank. A convenient, fast and safe soil amendment and fertilizer storage and dispensing system is thus desired."}
{"text": "In liquid storage tanks having floating roofs, the floating roof is constructed at a slightly smaller diameter than the side wall of the tank to define an annular space that permits some freedom of movement as the roof is raised and lowered by the fluid product stored in the tank. The annular space is sealed by a mechanical ring seal of which there are a number of designs such as those disclosed in Moyer, U.S. Pat. No. 2,696,930 and Lippiello et al., U.S. Pat. No. 5,103,992.\nSeals are necessary in most installations where the liquid contents vaporize and cause loss of product and pollution. Improved seals have been developed that permit the storage of liquids that are more volatile than were previously stored in floating roof tanks resulting in considerable savings over fixed roof designs.\nAfter extended periods of use, seals tend to break down and need to be replaced. In some older designs, certain components such as shoe plates were constructed from arcuate segments fastened together to form essentially one continuous ring around the tank. To replace or repair any ring seal component, the tank had to be taken out of service because environmental concerns prohibited use of a tank containing harmful vapor without a ring seal. Taking a tank out of service is time-consuming, hazardous and expensive."}
{"text": "The central nervous system (CNS) is composed of the brain contained in the cranium, and the medulla spinalis or spinal cord in the vertebral canal. The brain contains 12 cranial nerves. The spinal nerves spring from the medulla spinalis, and are transmitted through the intervertebral foramina. The spinal cord, which originates immediately below the brain stem, extends to the first lumbar vertebra, designated L1. Beyond L1 the spinal cord becomes the cauda equina. The spinal cord is comprised of 31 pairs of nerves comprising 8 pairs of cervical spinal nerves, 12 pairs of thoracic spinal nerves, 5 pairs of lumbar spinal nerves, 5 pairs of sacral spinal nerves, and 1 pair of coccygeal spinal nerves. The first cervical nerve (called the suboccipital nerve) emerges from the vertebral canal between the occipital bone and the atlas; the eighth issues between the seventh cervical and first thoracic vertebra. Each nerve is attached to the medulla spinalis by two roots, an anterior or ventral, and a posterior or dorsal, the latter being characterized by the presence of a ganglion, the spinal ganglion. Each spinal nerve is formed from a ventral or anterior root comprising axons from motor neurons and a dorsal or posterior root comprising nerve fibers from sensory neurons. The anterior root (radix anterior; ventral root) emerges from the anterior surface of the medulla spinalis as a number of rootlets or filaments (fila radicularia), which coalesce to form two bundles near the intervertebral foramen. The posterior root (radix posterior; dorsal root) is larger than the anterior due to the greater size and number of its rootlets which are attached along the posterolateral furrow of the medulla spinalis and unite to form two bundles which join the spinal ganglion. The posterior root of the first cervical nerve is smaller than the anterior. The spinal ganglia (ganglion spinale) are collections of nerve cells on the posterior roots of the spinal nerves. Each ganglion is oval in shape, reddish in color, and its size is in proportion to that of the nerve root on which it is situated. It is bifid medially where it is joined by the two bundles of the posterior nerve root. The ganglia are usually but not always found in the intervertebral foramina, immediately outside the points where the nerve roots perforate the dura mater. The ganglia of the first and second cervical nerves lie on the vertebral arches of the atlas and axis respectively, the ganglia of the sacral nerves are inside the vertebral canal, while the ganglion on the posterior root of the coccygeal nerve lies within the sheath of dura mater. When a dorsal root and a ventral root unite, they form a spinal nerve providing both motor and sensory utility. Spinal nerves split into two main branches, the dorsal ramus which innervates the skin of the back and deep back muscles, and the ventral ramus which innervates everything else from the neck inferiorly and also forms nerve plexuses which are a network of converging and/or diverging nerve fibers. The ganglia are comprised of unipolar nerve cells, and from these the fibers of the posterior root originate. Two other forms of cells are also present, i.e., the cells of Dogiel, whose axons ramify close to the cell (type II, of Golgi), and which are distributed entirely within the ganglion; and multipolar cells similar to those found in the sympathetic ganglia. The ganglia of the first cervical nerve may be absent, while small aberrant ganglia consisting of groups of nerve cells are sometimes found on the posterior roots between the spinal ganglia and the medulla spinalis. Each nerve root has a covering comprising pia mater, and is loosely invested by the arachnoid, the latter being prolonged as far as the points where the roots pierce the dura mater. The two roots pierce the dura mater separately, each receiving a sheath from this membrane; where the roots join to form the spinal nerve this sheath is continuous with the epineurium of the nerve.\nThe largest nerve roots and the largest spinal nerves are those of the lower lumbar and upper sacral nerves, and are attached to the cervical and lumbar swellings of the medulla spinalis.\nThese nerves are distributed to the upper and lower limbs. Their individual filaments are the most numerous of all the spinal nerves. The roots of the coccygeal nerve are the smallest. Immediately beyond the spinal ganglion, the anterior and posterior nerve roots unite to form the spinal nerve which emerges through the intervertebral foramen. Each spinal nerve receives a branch (gray ramus communicans) from the adjacent ganglion of the sympathetic trunk, while the thoracic and the first and second lumbar nerves each contribute a branch (white ramus communicans) to the adjoining sympathetic ganglion. The second, third, and fourth sacral nerves also supply white rami. These are not connected with the ganglia of the sympathetic trunk, but run directly into the pelvic plexuses of the sympathetic nerve system. A typical spinal nerve contains fibers belonging to two systems, i.e., the somatic nerve system, and the sympathetic or splanchnic nerve system, together with nerve fibers connecting these systems with each other.\nThe somatic fibers are efferent and afferent. The efferent fibers originate in the cells of the anterior column of the medulla spinalis, and run outward through the anterior nerve roots to the spinal nerve. They convey impulses to the voluntary muscles, and are continuous from their origin to their peripheral distribution. The afferent fibers convey impressions inward, for example from the skin, and originate in the unipolar nerve cells of the spinal ganglia. The single processes of these cells divide into peripheral and central fibers, and the latter enter the medulla spinalis through the posterior nerve roots.\nThe sympathetic fibers are also efferent and afferent. The efferent fibers, preganglionic fibers, originate in the lateral column of the medulla spinalis, and are conveyed through the anterior nerve root and the white ramus communicans to the corresponding ganglion of the sympathetic trunk where they may end by forming synapses around its cells, or may run through the ganglion to end in another of the ganglia of the sympathetic trunk, or in a more distally placed ganglion in one of the sympathetic plexuses. They end by forming synapses around other nerve cells. From the cells of the ganglia of the sympathetic trunk other fibers, postganglionic fibers, take origin; some of these run through the gray rami communicantes to join the spinal nerves, along which they are carried to the blood vessels of the trunk and limbs, while others pass to the viscera, either directly or after interruption in one of the distal ganglia. The afferent fibers are derived partly from the unipolar cells and partly from the multipolar cells of the spinal ganglia. Their peripheral processes are carried through the white rami communicantes, and after passing through one or more sympathetic ganglia without interruption end in the tissues of the viscera. The central processes of the unipolar cells enter the medulla spinalis through the posterior nerve root and form synapses around either somatic or sympathetic efferent neurons to complete reflex arcs. The dendrites of the multipolar nerve cells form synapses around the cells of type II (cells of Dogiel) in the spinal ganglia. By this path an original impulse is transferred from the sympathetic to the somatic system, through which it is conveyed to the sensorium.\nAfter emerging from the intervertebral foramen, each spinal nerve gives off a small meningeal branch which reenters the vertebral canal through the intervertebral foramen and supplies the vertebra and their ligaments as well as the blood vessels of the medulla spinalis and its membranes. The spinal nerve then splits into a posterior or dorsal, and an anterior or ventral division, each receiving fibers from both nerve roots.\nThe spinal cord provides a means of motor and sensory communication between the brain and the nerves of the peripheral nervous system (PNS) which includes the somatic nervous system (SNS) and the autonomic nervous system (ANS). The somatic nervous system is voluntary, includes the nerves serving the musculoskeletal system and the skin, and reacts to outside stimuli affecting the body. The autonomic nervous system is involuntary, and controls and maintains homeostasis or normal function. The autonomic nervous system is comprised of the sympathetic nervous system associated with the flight or fight response, and the parasympathetic nervous system responsible for maintaining regular life functions such as heartbeat and breathing during normal activity. Nerve roots pass out of the spinal canal through the intervertebral foramen and distribute to the body anteriorly for motor activity or posteriorly for sensory activity. The anterior nerve divisions supply the front of the spine including the limbs while the posterior nerve divisions are distributed to the muscles behind the spine.\nCerebrospinal fluid is secreted from the choroids plexus in the brain and is present in the brain ventricles, in the spinal canal, and in the spinal cord. The fluid circulates among these tissues and cushions the tissues in the CNS from the effects of traumatic injury. The CNS in a normal adult contains about 150 milliliters of cerebrospinal fluid.\nThe brain and spinal cord are covered and protected by meninges membranes comprising strong connective dura mater tissue, or dura, as a gray outer layer of the spinal cord and nerve roots; arachnoid mater, thinner than the dura mater; and pia mater, the innermost layer covering the nerves as a delicate and highly vascular membrane providing blood to the neural structures.\nThe dura membrane that covers the spine and nerve roots in the neck is surrounded by the epidural space. Nerves pass through the epidural space to the neck, shoulder and arms. Inflammation of these nerve roots may cause pain in these regions due to irritation from a damaged disc or from contract with the bony structure of the spine.\nInjury to the nerves in the central nervous system such as nerves in the spinal cord can result in functional impairment which can take the form of permanent loss of sensation and paraplegia. Most of the deficits associated with spinal cord injury such as traumatic spinal cord injury, crush injury, or lesion in the spinal nerves result from cell death and the loss of axons in the spinal neuronal population that are damaged in the central nervous system (CNS) which is comprised of nerves in the spinal cord and brain. Axons do not otherwise regrow across a lesion site in a damaged nerve in the CNS. Neurodegenerative diseases of the CNS are also associated with cell death and axonal loss. Representative diseases of the CNS include those that can result in impairment include stroke, human immunodeficiency virus (HIV) dementia, prion diseases, Parkinson's disease, Alzheimer's disease, multiple sclerosis, traumatic brain injury, and glaucoma. The ability to stimulate growth of axons from the injured, damaged, diseased or otherwise affected neuronal population would improve recovery of lost neurological functions, and protection from cell death can limit the extent of damage in the CNS. For example, following a white matter stroke, axons are damaged and lost, even though the neuronal cell bodies are alive, and stroke in gray matter kills many neurons and non-neuronal (glial) cells.\nEffective neuroprotective and neuroregenerative agents are desirable to potentially limit damage and to induce repair to the CNS. Compounds which promote axon growth in nerve cells of the CNS are especially desirable for treatment of damaged, injured, and/or diseased nerves in the CNS. Compounds which promote axon growth in nerve cells of the peripheral nervous system (the PNS) are also especially desirable for treatment of damaged, injured, and/or diseased nerves in the PNS.\nTraumatic injury of the spinal cord can result in permanent functional impairment. Axon regeneration in nerve cells of the CNS does not occur at the site of injury or lesion site in the mammalian CNS because growth inhibitory proteins present in anatomical structures immediately proximal to nerves as substrate-bound proteins block axon growth. Anatomical structures such as myelin sheaths immediately proximal to nerves are referred to as inhibitory substrates. While compounds such as trophic factors can enhance neuronal differentiation and stimulate axon growth in tissue culture, most factors and compounds that enhance growth and differentiation cannot promote axon regenerative growth on inhibitory substrates.\nA compound which stimulates axon growth in tissue cell culture and induces axon growth on growth inhibitory substrates is potentially useful for therapeutic use in axon regeneration when applied to cells residing in the CNS, such as directly to the site of a lesion in the CNS. Trophic and differentiation factors that stimulate growth on permissive substrates, that is, in the absence of inhibitory proteins or in the absence of inhibitory substrates, in tissue culture include neurotrophins such as nerve growth factor (NGF) and brain-derived growth factor (BDNF). Neither NGF nor BDNF promotes neurite growth or axon regeneration in nerve cells on inhibitory substrates, and neither is effective in promoting axon regeneration in nerve cells in the CNS in vivo.\nCell death can occur by two major mechanisms, necrosis and apoptosis. While necrotic cell death results in cell lysis and release of cell contents, cellular apoptosis is programmed cell death that results in the relatively tidy packaging of cells which die with the prevention of release of cellular contents. Apoptosis is characterized morphologically by cell shrinkage, nuclear pyknosis, chromatin condensation, and blebbing of the plasma membrane. Traumatic injury and ischemia can lead to apoptosis of both neurons and non-neuronal cells, and this cell death is responsible for functional deficits after injury or ischemia. A cascade of molecular and biochemical events is associated with apoptosis including activation of an endogenous endonuclease that cleaves DNA into oligonucleosomes detectable as a ladder of DNA fragments in agarose gels. Apoptotic endonucleases not only affect cellular DNA by producing the classical DNA ladder but also generate free 3′-OH groups at the deoxyribose ends of these DNA fragments. A technique called Tunel labeling labels DNA fragments as a means to detect apoptotic cells.\nRho kinase, an enzyme that resides in the interior of cells such as nerve cells, is a target for treatment of cancer, metastasis, and hypertension. Rho kinase inhibitors may be useful to treat eye diseases such as glaucoma, as well as to inhibit cancer cell migration and metastasis.\nRho kinase antagonists may be useful in treatment of hypertension, asthma, and vascular disease such as thrombosis.\nThe superfamily of small GTP-binding proteins or small G-proteins can be divided according to the similarity of amino acid sequences into 5 groups of Ras, Rho (short for Ras homologue), Rab, Arf and others. The small GTP-binding proteins have a molecular weight of 20,000-30,000 Daltons, specifically bind GDP and GTP, and exhibit GTPase activity by hydrolyzing bound GTP. Rho is specifically ADP ribosylated and inactivated by C3, a botulinum toxin.\nProtein kinases are proteins that phosphorylate and control the activity of other proteins by transferring a phosphate from ATP (adenosine triphosphate) to an amino acid, e.g. serine, threonine, or tyrosine, on another (target) protein. Target proteins regulated by phosphorylation include enzymes that transduce signals in a cellular environment and enzymes that turn certain genes on or off and can thereby regulate progression of many different diseases.\nRho kinase regulates cell cytoskeleton organization, cell adhesion and cell motility, and a cell's cycle. Rho kinase is a serine/threonine kinase, and is a Rho binding protein or an effector of Rho which is a GTPase, which catalyzes the reaction: GTP (Guanosine 5′-triphosphate)+H2O (water) to GDP (Guanosine 5′-diphosphate)+phosphate ion, and is linked to a cell membrane in its active state. Inhibition of Rho kinase by a compound of this invention can have potential therapeutic applications in a mammal, such as reduction of tumor metastasis in cancer, relaxation of vascular tension in cardiovascular disease, and reduction of ocular pressure in glaucoma, among other applications. Two different Rho kinase inhibitors, Fasudil™ and Radicut™, are in use in humans for treatment of stroke. Rho kinase is called ROK (Rho-associated kinase) and ROCK (Rho associated Rho kinase), ROKa, or ROCKI. Isoforms of ROCK exist: ROCKII (has 64% sequence identity with ROCKI, and ROKβ is 90% identical. The protein has three important domains: a Rho-binding domain (RB), a C-terminal PH (Pleckstrin homology) domain, and a catalytic domain.\nROCKI and ROCKII are each activated by Rho. An important distinction between ROCKI and ROCKII comprises their respective in vivo tissue distributions in a mammal (e.g., see Nakagawa et al., FEBS Lett. 1996 Aug. 26; 392(2):189-93). Analysis of relevant mRNA concentration levels has shown that ROCKI has widespread tissue distribution, but that there is relatively little ROCKI in brain and in skeletal muscle. Expression of ROCKII is high in brain, heart and lung, but is relatively low in liver, stomach, spleen, kidney and testis. Western blots used to study protein concentration levels show ROCKII concentration levels to be low in liver and kidney (Wibberley et al. 2003, British J. Pharmacology 138:757-766), whereas ROCKII is highly expressed in brain as compared to ROCKI (Leemhuis et al. 2002, J Pharmacol Exp Ther. 300:1000). Therefore, for therapeutic and diagnostic use in the CNS of a mammal, an inhibitor that is specific for inhibition of ROCKII would be desirable.\nIn this regard, known Rho kinase inhibitor Y27632 [(R)-(+)-trans-4-(1-aminoethyl)-N-(4-pyridyl)cyclohexane-carboxamide as a dihydrochloride salt] inactivates (inhibits) both ROCKI and ROCKII (Ishizaki et al. 2000, Molecular Pharmacology 57:976).\nIn addition, a Rho kinase inhibitor that is in clinical use for treatment of stroke, Radicut™ (edaravone; 3-methyl-1-phenyl-2-pyrazoline-5-one; 5-methyl-2-phenyl-2,4-dihydro-3H-pyrazol-3-one), is known to have kidney toxicity. An estimated 180,000 patients have taken the drug Radicut™ since it was approved for use in 2001 in Japan. Ninety-three were later struck down with kidney failure and 40 died. A Rho kinase inhibitor specific for ROCKII may not be expected to show this adverse effect profile in part because of the low abundance of ROCKII in liver and kidney (Wibberley et al. 2003, British J. Pharmacology 138:757-766).\nNearly all protein kinase inhibitors that have been developed are ATP-competitive, and for this reason the “drug concentration required for 50% inhibition” (IC50) of a protein kinase depends on the concentration of ATP used in the assays. A drug concentration required to suppress phosphorylation of a target substrate in a cell can be higher than the drug concentration required for inhibition of a protein kinase in vitro.\nA compound known to inhibit the activity of Rho kinases is (R)-trans-4-(ethan-1′-amino)-N-(4″-pyridyl)cyclohexane carboxamide dihydrochloride known as Y-27632, which is available from Sigma Chemical Company and from Calbiochem in powder form and which can be dissolved in DMSO (dimethyl sulfoxide, Sigma) and which can be useful as a reference Rho kinase inhibiting compound.\nRho kinase regulates axon growth and regeneration in nerve cells, cell motility and metastasis, smooth muscle contraction, and apoptosis, and is a target for therapeutic treatment in many disease applications, including repair in the central nervous system. Rho kinase is activated by Rho and Rho kinase inhibitors block Rho signaling. Mutations of Rho family regulatory proteins have been found in clinical oncology samples which suggests that perturbation or alteration or interruption of, or interference with, the Rho signaling pathway can be a useful therapeutic modality. Examples with specificity for Rho include the DLC1 gene in hepatocellular carcinoma, p-190-A, which is in a region that is altered in gliomas and astrocytomas, GRAF, which has loss of function mutations in leukemia, and LARG, which is found in some gene fusions found in acute myeloid leukemia. Genetically engineered point mutations can activate RhoA and induce cellular transformation in vitro.\nRho kinase inhibitors have widespread potential for use in the treatment of neurodegenerative diseases, particularly if the Rho kinase inhibitors have the property to enhance plasticity and axon regeneration in neurons. However, there is much scientific evidence for a direct link between Rho signaling and neurodegenerative disease. In an animal model of Alzheimer's Disease (AD), there is clear evidence that Rho kinase inhibitors can reduce the pathological hallmarks of the disease.\nA trans-4-amino(alkyl)-1-pyridylcarbamoylcyclohexane, designated as Y27632, is a Rho kinase inhibitor and available from Calbiochem. Y27632 is described in U.S. Pat. No. 4,997,834, the entire content of which is incorporated herein by reference. Y27632 has been used to demonstrate that inhibition of Rho kinase is effective in preventing metastasis. Other compounds are described in U.S. Pat. No. 5,478,838, the entire content of which is incorporated herein by reference.\nRho signaling antagonists can be useful in treatment of hypertension. For example, Y-27632 can relax smooth muscle and increase vascular blood flow. Y-27632 is a small molecule that can enter cells, and is not toxic in rats after oral administration of 30 mg/kg for 10 days. Y-27632 reduces blood pressure in hypertensive rats, but does not affect blood pressure in normal rats.\nA number of Rho kinase inhibitors are known. For example, the compound NHM-1152 can act as vascular relaxant in vascular vasospasm. The compound known as hydroxy fasudil may find use in the treatment of stroke after intravenous application and can reduce infarct volume, improve outcomes, can have anti-ischemic properties in vasospastic angina, and inhibits neutrophil migration in ischemic brain. A fasudil compound called HA-1077 is an antivasospasm agent which improves cerebral hemodynamic activity, inhibits production of superoxide anion by neurotrophils, and may find use in the treatment of spinal cord injury, stroke, subarachnoid hemorrhage, and cerebral infarction. U.S. Pat. No. 6,218,410, the disclosure of which is hereby incorporated by reference in its entirety, describes a Rho kinase inhibitor. U.S. patent application Ser. No. 10/022,301 (McKerracher et al) published as U.S. 2002/0119140, the disclosure of which is hereby incorporated by reference in its entirety, describes Rho family antagonists and their use to block inhibition of neurite outgrowth. Canadian Patent applications 2,304,981 (McKerracher et al) and 2,325,842 (McKerracher), the disclosure of each of which is hereby incorporated by reference in its entirety, disclose the use of Rho antagonists such as for example C3 and chimeric C3 proteins as well as substances selected from among known trans-4-amino(alkyl)-1-pyridyl-carbamoylcyclohexane compounds or Rho kinase inhibitors for use in the regeneration of axons. C3 inactivates Rho by ADP-ribosylation and can be relatively non-toxic to nerve cells at therapeutically effective doses. U.S. Pat. Nos. 4,857,301, and 5,741,792, the disclosure of each of which is hereby incorporated by reference in its entirety, describes a number of 4-substituted piperidine compounds. U.S. Pat. No. 6,140,333, the disclosure of which is hereby incorporated by reference in its entirety, describes certain piperidine compounds, one of which is a 4-aminomethylpiperidine acylated at the piperidine ring nitrogen by a group containing a carbon adjacent to the carbonyl group to which carbon is attached an aryl group and a hydroxyl group. U.S. Pat. No. 6,020,352, the disclosure of which is hereby incorporated by reference in its entirety, describes the use of certain 1-phenyl-2-piperidinoalkanol derivatives to treat ischemic disorders of the retina and optic nerve.\nU.S. Pat. Nos. 4,849,521, 4,584,303, 4,866,077, 4,933,353, and 6,169,097, the disclosure of each of which is hereby incorporated by reference in its entirety, disclose methods to prepare a number of piperidines with substituents attached to one or more of the ring carbon atoms of a piperidine. U.S. Pat. No. 6,545,022, the disclosure of which is hereby incorporated by reference in its entirety, describes a method to prepare certain 4-substituted-4-aminoalkylpiperidines including certain 4-substitutedmethylene-4-aminomethylpiperidines."}
{"text": "In recent years, the availability of advanced location-specific agricultural application and measurement systems (used in so-called “precision farming” practices) has increased grower interest in determining spatial variations in soil properties and in varying input application variables (e.g., planting depth) in light of such variations. However, the available mechanisms for measuring properties such as temperature are either not effectively locally made throughout the field or are not made at the same time as an input (e.g. planting) operation.\nThus, there is a need in the art for a method for monitoring soil properties during an agricultural input application."}
{"text": "The present invention relates generally to drilling a borehole, and more particularly to a method for real time geostreering the drilling of a borehole.\nIn the search for oil and gas producing formations to meet the ever increasing demand for energy, the oil industry constantly strives to develop techniques to maximize the recovery from oil and gas producing reservoirs. Oftentimes, the recoverable oil and gas in the reservoirs is in thin oil and gas bearing formations or beds, for example, three to twenty feet thick, that requires the drilling of horizontal wells.\nTechniques have been developed to steer the drill bit to maintain the borehole within a target formation. For example, various steering techniques use information from offset wells to compare with logging data obtained while drilling the well. A logging while drilling (LWD) tool provides real time borehole measurements while drilling through a formation. The LWD measurements are used to make drill bit trajectory adjustments to maintain the borehole in the target formation. Use of real time information in this manner is generally referred to in the oil industry as geosteering."}
{"text": "Approaches described in this Background section may include approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section.\nTargeted marketing techniques assist a variety of entities, such as product manufacturers, service providers, and retailers, in efficiently persuading consumers to purchase certain products and service. A marketing entity may include, without limitation, product manufacturers, service providers, retailers, third-party marketing firms acting on behalf of another entity, non-profit organizations, and so forth. The marketing entity may identify characteristics of people to which the entity wishes to “market” a product or service. For example, an entity that sells product X may determine that it wants to market to persons that have purchased a complimentary product Y within a timeframe of Z. Of course, one problem with targeted marketing is that it is sometimes difficult to ascertain exactly what characteristics to target.\nOnce targeted characteristics have been identified, the entity may then attempt to “market” to persons who have those characteristics. Unfortunately, it is also generally difficult to ascertain whether a given person actually has the targeted characteristics and/or how to actually reach specific persons that have those specific characteristics. For example, a particular retailer may have no idea that a potential consumer recently purchased product Y from another store, and thus be unable to specifically target marketing for Product X to that potential consumer. Large-scale targeted marketing is therefore typically difficult to achieve with a high degree of efficiency. Rather, a marketing entity is typically limited to less efficient targeting techniques, such as targeting specific zip codes or web sites, thus resulting in the marketing entity wasting marketing resources on many persons who do not have the targeted characteristics.\nOne common marketing technique involves incentivizing costumer behavior through targeted promotional offers. An “offer” may be construed to be a promise by an offer provider, which may be any entity engaged in targeted marketing, to provide a consumer with a benefit under a certain set of implicit or explicit conditions known as terms. Example benefits include, without limitation, a monetary gift, a discount applied to the purchase of one or more products or services, free products and/or services, access to additional offers, and so forth. Example terms may include any one or more of, without limitation, actions that the consumer must perform, a specific set of item(s) that must be purchased by the consumer, a specific amount of money that must be spent by the consumer, a timeframe during which the offer is valid, and so forth. A customer benefit of an offer is realized during a process sometimes called “redemption,” which typically takes place during a transaction between the consumer and a retailer. When the consumer wishes to engage in a transaction the retailer determines whether the conditions of the offer have been met. If so, the retailer, which may be any merchant or other entity that sells products or services, provides the benefit to the consumer on behalf of the offer provider.\nIn some embodiments, the retailer is in fact the offer provider. In other embodiments, the retailer may then seek compensation for providing the benefit from the offer provider by means of a process known as clearing. During the clearing process, the retailer provides evidence to a third-party clearinghouse, or directly to the offer provider, that the retailer provided the benefit. Assuming the evidence is sufficient, the third-party clearinghouse, or offer provider, provides compensation to the retailer. In embodiments involving a third-party clearinghouse, the offer provider may provide the third-party clearinghouse with funds in advance. The offer provider also sends to the third-party clearinghouse funds from which the compensation is provided, or reimburses the third-party clearinghouse for providing the compensation.\nA frequent implicit or explicit term of an offer is that the retailer must be presented with a specific coupon during a transaction. In general, a coupon is a certificate, document, or other physical or electronic construct that entitles a consumer to accept an offer described or referenced by the coupon, thereby realizing the benefit of the offer. A coupon sometimes takes a “hard copy” form, such as a paper certificate, with printed images and/or text describing terms of the offer. The process of the consumer accepting an offer by presenting, referencing, or otherwise providing a coupon while purchasing, contracting, or otherwise transacting with another party involves “redeeming” the coupon. For example, a consumer may redeem a hard copy of a coupon by handing the copy to a clerk during a purchase at a retail store. The clerk may then provide the consumer with the subject of the offer, such as a discounted price, item, service or gift.\nOne technique for distributing coupons is to include printed coupons with newspapers, magazines, or other items that are distributed to consumers. One example of an item with which coupons are distributed is a printed receipt. For example, some retailers print receipts at a point of sale on register paper on which coupons have been pre-printed. As another example, some retailers print coupon(s) on a receipt at the time of the transaction for which the receipt is printed, thereby allowing the retailers to dynamically select which coupon(s) appear on the receipt based on the product(s) that were purchased during the transaction.\nRecent distribution techniques now provide consumers with opportunities to print their own coupons. For example, a number of websites provide search engines or catalogs with which consumers may locate offers and then print coupons for the offers they find. The printed coupons may be used in the same manner as any other coupon.\nOther recent distribution techniques involve creating digital coupons. One such technique involves creating unique digital coupons that are saved to an account associated with the consumer, such as a store loyalty account. The consumer may redeem such digital coupons during online or physical transactions by presenting an account identifier, such as a store loyalty card or an oral identification of the consumer's telephone number, for the associated account. Since many consumer accounts are tied to card-based identifiers, such as store loyalty cards or credit cards, the process of storing a digital coupon identifier to an account may also be referred to as saving a coupon to a card. Some examples of such techniques are described in U.S. application Ser. No. 12/878,231, filed Sep. 9, 2010, the entire contents of which are hereby incorporated by reference for all purposes as if fully set forth herein.\nAnother digital coupon-based technique involves creating unique digital coupons that may be stored on a computing device. The digital coupons may be transmitted from the computing device to a point-of-sale during a transaction using any of a variety of mechanisms. For example, information about the digital coupon may be uploaded to the point-of-sale during an online transaction involving the computing device. As another example, information about the digital coupon may be transmitted wirelessly from a smartphone to a receiving component coupled to a checkout register during a transaction at a brick-and-mortar store.\nConsidering the current state of the industry and market for digital offers (including digital coupons and other digital promotional discount vehicles), there is a need for improved systems, methods and technology platforms for generating, processing, storing, managing and delivering digital offers, for enhancing and customizing the digital offer experience of consumers in general, for increasing the value and benefits derived by consumers from digital offers, for improving the ability of offer providers, offer distributors and retailers to customize digital offer selection, delivery, utilization, management, monetization and redemption analysis, and for otherwise increasing the efficiency and financial return of the digital offer industry and market."}
{"text": "Conventionally, as a foam polyurethane material, for example, Patent Document 1 below has proposed a microcellular polyurethane elastomer produced by allowing a polyisocyanate compound to react with polyol, having an overall density of 100 kg/m3 or more, and having an overall density, a compression set, and a cell diameter on skin surface in specific relation.\nFurthermore, as the foam polyurethane material, for example, Patent Document 2 below has proposed a polyurethane foam produced from materials including polyisocyanate including bis(isocyanatomethyl) cyclohexane, polyol, a blowing agent, and a catalyst."}
{"text": "1. Field of the Invention\nThe present invention relates to a three-dimensional-position measuring method and a three-dimensional-position measuring apparatus.\nThe present invention also relates to a three-dimensional shape measuring method and a three-dimentional-shape measuring apparatus for measuring the position of an object and a three-dimensional-shape of the object using visual information of the three-dimensional object.\nThe present invention also relates to an image pickup (or picture taking) apparatus and a method for synthesizing respective picked-up images (or taken pictures) into one image (or picture), the respective picked-up images (or taken pictures) having been picked up (or taken) in a manner in which a portion of each image corresponds to a portion of another image, these images (or pictures) being images (pictures) present on an object plane to be input. Such an apparatus and a method enable picking up an image (or taking a picture) in a case where the object plane, an image (or picture) of which is to be input, is large or in a case where an object having the object plane, an image (or picture) of which is to be input, cannot be moved.\n2. Descriptions of the Related Art\nIn various fields such as surveying, production of a CAD model and so forth, measurement of a three-dimensional shape of an object has been demanded. As a method for measuring a three-dimensional shape, there is an active method in which a reference pattern or the like is projected on an object, and an image of distortion of the pattern is picked up. However, according to such a method, the object may be affected by projection of the reference pattern, and the effect of the reference pattern varies in accordance with the environment. Because such demerits exist, application of such a method is significantly limited. A passive method, by which an object is not affected, for measuring the shape of an object has been also used. Such a method for measuring a three-dimensional shape is a method in which, basically, corresponding points are extracted from images which have been picked up from two different points of view. Then, the three-dimensional position of the corresponding points is calculated in accordance with the principle of triangulation. In this method, two image pickup apparatuses are fixed and the images of an object are obtained from two different points of view. In another method, one image pickup apparatus is moved on rails a predetermined distance and images of an object are obtained from two different points of view. In these two methods, the relationship of the positions of the two different points of view and the change of the orientation of the image pickup apparatus(es) are treated as already-known information. Instead, there is another method in which an image pickup apparatus is moved in a condition where the apparatus is held by hands of a person, and the relationship of the positions of two different points of view and the change of the orientation of the image pickup apparatus(es) are treated as not-already-known information.\nIn methods such as those described above in which the relationship of the positions of two different points of view and the change of the orientation of image pickup apparatus(es) are treated as already-known information, a measuring apparatus therefor should be a large-scale apparatus. Therefore, application of such a method is limited. In the method in which the relationship of the positions of two different points of view and the change of the orientation of image pickup apparatus(es) are treated as not-already-known information, the method can be applied to various uses because an object is not affected by such a method and a measuring apparatus therefor can be a small-scale apparatus. As disclosed in, for example, Japanese Laid-Open Patent Application Nos.5-196437, 6-129833, 7-181024, 7-244735 and so forth, various measuring methods and measuring apparatuses have been developed.\nIn the method disclosed in Japanese Laid-Open Patent Application No.5-196437, an image of a measurement point on an object is picked up by a camera in orthogonal projection, the orientation of the camera at the current time is obtained from a three-axis gyro fixed to the camera, and the thus-obtained information is supplied to a three-dimensional-information measuring apparatus. Thus, the three-dimensional information of the object is measured. In the method disclosed in Japanese Laid-Open Patent Application No.6-129833, using moving image information of an object, images of which have been picked up by a monocular camera, the motion of the camera is obtained, and the three-dimensional shape of the object is calculated. In the method disclosed in Japanese Laid-Open Patent Application No.7-181024, an optical flow, which indicates movement between images, is obtained from a plurality of images of an object, picked up by image pickup means in a manner in which the points of view of the image pickup means are different. Then, based on the optical flow, distance information is calculated. Further, in this method, the distance information is directly calculated by using a formula for a relationship between the time difference and the space gradient of the images and by using the movement information of the image pickup means itself. Then, these ways for obtaining the distance information are combined, and the three-dimensional shape is measured. The apparatus disclosed in the Japanese Laid-Open Patent Application 7-244735 obtains three-dimensional information as a result of temporally combining distance information obtained by using moving image information and movement of a camera.\nHowever, in the case where an integral value of output from an acceleration sensor or a gyro is used for obtaining position change or orientation change as disclosed in Japanese Laid-Open Patent Application Nos.5-196437 and 7-181024, because of use of the integral value, reliability of the information of the thus-obtained position change or orientation change is low. In the case where only image information is used as disclosed in Japanese Laid-Open Patent Application No.6-129833, it is difficult to extract precise corresponding points. In the case where, as disclosed in Japanese Laid-Open Patent Application No.7-244735, the movement speed of the image pickup apparatus is sensed, and the thus-obtained information is used in order to improve the measurement accuracy, the orientation change of the image pickup apparatus is not sensed. Because information concerning the orientation change is important for extracting corresponding points and shape measurement, significant improvement of the measurement accuracy cannot be expected.\nAn apparatus for inputting and restoring characters or a photograph present on a plane simply and with high accuracy is demanded. When an image of an ordinary size such as an image drawn on a paper sheet of A4 size, A3 size or the like is input, inputting the image with high accuracy can be performed by a copying machine, a scanner or the like. However, a copying machine, a scanner and so forth cannot be used for inputting, with high accuracy, a larger image, an image drawn on a medium which cannot be moved or the like, such as information drawn on a larger paper sheet, information drawn on a wall, information drawn on a panel or the like. As apparatuses and a method for inputting a larger image, an image drawn on a medium which cannot be moved, or the like, using an electronic camera, an electronic camera disclosed in Japanese Laid-Open Patent Application No.6-141228 and a panorama image synthesizing apparatus and a method for producing a panorama image disclosed in Japanese Laid-Open Patent Application No.9-90530 have been developed.\nThe electronic camera disclosed in Japanese Laid-Open Patent Application No.6-141228 forms the image of an object on an image pickup element via an image pickup lens and a mirror, picks up an image of the object intermittently several times in a manner in which a portion of a previously picked-up image is the same as a portion of a currently picked-up image. The electronic camera moves in a manner in which the currently picked-up image is appropriately connected with the previously picked-up image. In the panorama image synthesizing apparatus and the method for producing a panorama image disclosed in Japanese Laid-Open Patent Application No.9-90530, parameters for synthesizing images are calculated, and the images are synthesized based on the calculated parameters. Thereby, it is attempted to obtain a synthesized image with high accuracy using separate images picked up by the electronic camera.\nHowever, in the electronic camera disclosed in Japanese Laid-Open Patent Application No.6-141228, the mirror and a mechanism for driving the mirror are needed. Thereby, the scale of an apparatus therefor is large. Accordingly, it is difficult to install the apparatus in an electronic camera of an ordinary size.\nFurther, when an image on a paper sheet or a panel is picked up, the distance between a camera and an object is short. Therefore, in the panorama image synthesizing apparatus and the method for producing a panorama image disclosed in Japanese Laid-Open Patent Application No.9-90530, due to influence of perspective projection due to translation operation of the camera, distortion of the synthesized image occurs. Further, an image rotates because a person holds the electronic camera, and the image rotation causes distortion of the synthesized image."}
{"text": "Pencil sharpeners for sale online and at retail stores are typically versions that are either held in one's hand during use, or fixed in one place such as on a desk during use. U.S. Pat. No. 2,702,022 discloses both of these pencil sharpener versions.\nBoth pencil sharpener versions have drawbacks. The handheld version must be carried about by the user so it is available when needed for use. This also makes it vulnerable to being misplaced or lost. The desk-mounted version requires that users have to go to where it is located and possibly wait in line to use it.\nAccordingly, there remains a need in the art for an innovation that will overcome the deficiencies of past approaches and the problems that remain unsolved."}
{"text": "1. Field of the Invention\nThe present invention relates to a coupling structure of an optical connector, in particular, to that of an optical connector serving as an optical connecting device to be connected to a small-sized optical component which is, like an optical semiconductor module used in an optical communication system and others, fixed to a board.\n2. Description of Prior Art\nThere have been developed and used various types of optical connectors and some of them are disclosed in Japanese Patent Provisional Publications No. 211276 of 1997, No. 43453 of 1997, and No. 265026 of 1997. Among them, in an embodiment where an optical semiconductor module disclosed in Japanese Patent Provisional Publication No. 211276 of 1997 and an optical connecting device using the optical semiconductor module are employed, the embodiment has an optical semiconductor module, which has a package body accommodating a semiconductor laser and an optical fiber optically coupled thereto and double side walls respectively provided with a recess and formed in the traveling way of the laser light radiated from the semiconductor laser, and a simplified connector to be coupled with said module, having a recess in the inner wall holding a ferrule to which one end of the optical fiber is fixed, and being provided at its front end with a protuberance, and optical fibers are optically coupled with each other through the coupling of the ferrules by coupling the simplified connector with the optical semiconductor module.\nThough these conventional optical connectors are suitable, to some extent, for producing a large amount of low cost and small-sized optical semiconductor modules indispensable in optical fiber communication systems, they usually still have some following problems. That is to say, special tools are necessary for centering the ferrules during the assembly work of coupling the module with the connector, and this makes the assembly difficult. In case that the coupling sequence for a mate optical semiconductor module starts from a ferrule and there is no means to ensure the centering for the ferrules, a successful coupling is hard, the ferrule or a sleeve is sometimes damaged, and the coupling work eventually becomes difficult. Further, even if there is a guide means for centering the ferrules, this guide means does not sufficiently answer the need of low cost because the guide means is provided as a member completely independent of a lock means for locking the coupling of the mate optical semiconductor module, and this causes an increase in the number of components. Furthermore, the optical connection has been unstable due to the lack of means for ensuring the ferrule alignment after the coupling.\nAn object of the present invention is to provide a structure allowing the conventional technological problems described above to be eliminated for an optical semiconductor module, an optical connector, or a shape adapting optical connector."}
{"text": "The present invention relates to a pressurizing image forming apparatus which makes use of a photo- and pressure-sensitive heat-developable material having a substrate and a layer formed on the substrate from a substance which is photosensitive and developable by application of heat and capable of fixing the developed image by application of pressure. More particularly, the present invention is concerned with a pressurizing image forming apparatus which fixes, by application of pressure, an image which has been formed on the heat-developable material of the type mentioned above through exposure and developed by application of heat.\nAn image recording material employing microcapsules containing a photosensitive composition has been known. An example of such a recording material is disclosed in Japanese Patent Laid-open No. 179,836/1982. This recording material has a substrate which carries capsules made of a synthetic polymeric resin and containing a vinyl compound, a photopolymerization initiator and a coloring precursor.\nIn recording an image using this recording material, the material is first exposed so that the microcapsules are hardened in conformity with the form of the image. Then, as the recording material is pressed, the microcapsules which have not been hardened are ruptured so as to release the coloring agent precursor, thus forming a color image. This image recording system provides a high quality of the recorded image with a dry-type simple processing, but suffers from a disadvantage in that the photosensitivity is much smaller as compared with the system which makes use of silver halide.\nUnder this circumstance, the present applicant has proposed, in the specification of Japanese Patent Laid-open No. 275742/1986, a novel recording material which has a high photosensitivity and which ensures a high quality of the recorded image with a simple dry-type process. This recording material is a photo- and pressure-sensitive heat-developable material constituted by a carrier and a layer formed on the surface of the carrier, the layer containing at least a photosensitive silver halide, reducing agent, polymerizable compound and a color image forming substance, wherein the polymerizable compound and the color image forming substance are confined in common microcapsules.\nAn image recording method which makes use of this photo- and pressure-sensitive heat-developable material is disclosed in the specification of Japanese Patent Laid-open No. 278849/1986 filed by the same applicant. According to this method, the heat-developable material is first exposed so that a latent image is formed thereon in conformity with the image to be recorded. Then, the material is heated for development so that the polymerizable compound in the area where the latent image exists is polymerized to produce a polymeric compound, thus thermally setting the microcapsules. Then, the material is superposed on an image receiving material having an image receiving layer capable of receiving the color image forming substance, and is pressed to the image receiving material so that at least part of the microcapsules having no latent image is ruptured so as to transfer the color image forming substance to the image receiving material, thereby forming an image on the latter.\nIn order to obtain a clear image by the transfer of the heat-developed image from the photo- and pressure-sensitive heat-developable material to the image receiving material, it is essential that the heat-developable material be pressed uniformly onto the image receiving material under a given condition."}
{"text": "1. Technical Field\nThe present invention relates to technology of ejecting a liquid such as an ink.\n2. Related Art\nA liquid ejecting head that ejects a liquid such as an ink has been proposed in the related art. In a case where the liquid ejecting head is filled with the liquid (hereinafter, referred to as “initial filling”), air, with which a space inside the liquid ejecting head is filled, needs to be discharged. In addition, in a state in which the space inside the liquid ejecting head is filled with the liquid, it is important that bubbles mixed with the liquid be discharged. In this respect, JP-A-2002-144576 discloses a configuration in which a bubble outlet for discharging gases is disposed in a ceiling surface of a common liquid chamber in which a liquid that is supplied to a plurality of nozzles is stored.\nHowever, in the technology disclosed in JP-A-2002-144576, there is a possibility that the liquid stored in the common liquid chamber will leak through the bubble outlet or that bubbles from the bubble outlet will be mixed in the common liquid chamber."}
{"text": "Morphology of various organs of animals is constructed by epithelial tissues, and mesenchymal cells exist around the tissues. An epimorphin is a cell membrane protein that is expressed in mesenchymal cells, particularly in a high amount neighboring epithelial cells (Hirai, Y., et al., Cell, 69, pp. 471-481, 1992). It is considered that the progress of morphogenesis of the epithelial tissue requires a signal from the mesenchymal cells (Gumbiner, B. M., Cell., 69, pp. 385-387, 1992). Epimorphins have been cloned in humans, birds, and rats as well as in mice. The presences of isoforms are known which have a different sequence of a hydrophobic site (Zha, H., et al., Genomics, 37, pp. 386-389, 1996; Hirai, Y., et al., Biochem. Biophys. Res. Commun., 191, p. 1332-1337, 1993; Oka, Y., Developmental Biological Society, 1997, May).\nEpimorphin is known to be deeply involved in morphogenesis by epithelial tissues in mice, i.e., differentiation into hair on the fetal talon skin, and differentiation into luminal structure in the fetal lung (Hirai, Y., et al., Cell, 69, pp. 471-481, 1992; Koshida, S., et al., Biochem. Biphys. Res. Commun., 234, pp. 522-526, 1997). Moreover, epimorphin activates mesenchymal cells and promotes secretion of cytokines, IL-6 and IL-8 (Oka, Y., et al., Exp. Cell Res., 222, pp. 189-198, 1996). Recently, it was revealed that the addition of epimorphin to milk protein-producing SCp2 cells induces cell growth to form a branched duct structure (Hirai, Y., et al., J. cell Biol., 140, pp. 159-169, 1998). Epimorphin is expected to be effective in elucidating the mechanism of onset of diseases due to abnormality of organs, developing methods for diagnosing and treating said diseases, generating hairs, lumens, bones, and teeth, generating new blood vessels, and developing new methods for treatment of injuries (Zha, H., et al., Genomics, 37, pp. 386-389, 1996; Panaretto, B. A., Reprod. Fertil. Dev., 5, pp. 345-360, 1993; Matsuki, Y., et al, Archs. Oral Biol., 40, pp. 161-164, 1995)."}
{"text": "The present invention relates to charging condition control in a distance measuring apparatus for measuring the distance to an object and, more specifically, a method of controlling charging conditions suitably applied to an automatic focusing mechanism of a camera.\nConventionally, a distance measuring device which performs trigonometric calculation by projecting a light spot onto an object to be measured and receiving light reflected by the object using a position detection means such as a position sensitive detector (PSD) or the like is known. Further, another distance measuring device which circulates an accumulated charge using a ring-shaped charge transfer device, such as a CCD, to integrate reflected light of ON/OFF-projected light spots and skims a predetermined amount of charges of external light components other than the light spot has been proposed by Japanese Patent Publication No. 5-22843 and Japanese Patent Application Laid-Open No. 8-233571. The distance measuring device of this type can keep accumulating charges while circulating the accumulated charge if the level of the accumulated charge is not high enough, thereby it is possible to obtain signals of good S/N ratio.\nFIG. 12 is a diagram illustrating a configuration of a light-receiving unit 500 used in a distance measuring apparatus.\nNote, in FIG. 12, a photoelectric conversion (photo-receiving) device 520 of the light-receiving unit 500 is represented by three photoelectric conversion devices X, Y and Z, to simplify the explanation.\nAn image sensing apparatus having the light-receiving unit 500 operates in two different modes, namely, an active mode and a passive mode.\nThe active mode projects light onto an object 515 to be measured, the distance to which is to be measured, by turning on and off a light emit element (here, infrared light-emitting diode; IRED) 514 to emit light pulses, receives light reflected by the object 515, when the IRED 514 is on and when the IRED 514 is off, respectively, using the photoelectric conversion devices X, Y and Z, and stores the charges separately when the IRED 514 is on and when the IRED 514 is off. Whereas, the passive mode receives external light reflected by the object without turning on the IRED 514 using the photoelectric conversion devices X, Y and Z, and store the charges.\nThe distance measuring apparatus is of a hybrid-type capable of performing distance measuring operations both in the active mode and in the passive mode, and, when a reliable measurement result is not obtained in the active mode, then the distance is measured once again in the passive mode.\nFurther, the light-receiving unit 500 has a linear CCD 524 which includes ON-pixels 522x, 522y, and 522z and OFF-pixels 523x, 523y, and 523z, respectively corresponding to the photoelectric conversion devices X, Y and Z, a ring-shaped CCD 521 which includes a plurality of ON-pixels and OFF-pixels, and a skim CCD 523.\nThe charges obtained as a result of photoelectric conversion in the photoelectric conversion devices X, Y and Z are respectively transferred to the corresponding ON-pixels and OFF-pixels of the linear CCD 524 and stored. Thereafter, the charges are transferred to the ring-shaped CCD 521.\nNext, timing of charge transfer operation in the light-receiving unit 500 is explained with reference to FIG. 13.\nReferring to FIG. 13, the IRED 514 turns on and off in synchronization with the ON/OFF (High/Low) of a charging signal in the active mode, and the IRED 514 is kept off independent of the ON/OFF of the charging signal in the passive mode.\nBelow, the active mode is explained.\nFirst, charges obtained in the photoelectric conversion devices X, Y and Z while the charging signal is ON (i.e., High level) are transferred to the ON-pixels 522x, 522y, and 522z while an ON-pixel transfer signal is ON (i.e., High level).\nFurther, charges obtained in the photoelectric conversion devices X, Y and Z while the charging signal is OFF (i.e., Low level) are transferred to the OFF-pixels 523x, 523y, and 523z while an OFF-pixel transfer signal is ON (i.e., High level).\nIn this manner, charges due to projected light reflected by the object and external light are stored in the ON-pixels 522x, 522y, and 522z, while charges due to external light are stored in the OFF-pixels 523x, 523y, and 523z in the active mode.\nAfter the charges obtained in the photoelectric conversion devices X, Y and Z are transferred to the ON-pixels 522x, 522y, and 522z and the OFF-pixels 523x, 523y, and 523z, the charges are transferred to the ring-shaped CCD 521.\nTo transfer the charges to the ring-shaped CCD 521, a ring transfer signal is used. The ring transfer signal becomes High so that charges from the same pixel of the linear CCD 524 are always transferred to the same pixel of the ring-shaped CCD. Accordingly, charges outputted from the ON-pixel 522x, corresponding to the photoelectric conversion element X obtained during the charging signal is ON, for example, are accumulated.\nIn FIG. 13, the numerals 1, 2, 3, and so on, indicate the number of circulations. The number of circulations indicates the number of times charges are transferred to the ring-shaped CCD 521.\nMore specifically, in the first circulation, charges are transferred to the ring-shaped CCD 521 once, as shown in FIG. 14A, and the charges obtained in one charging operation are stored. In the second circulation, charges obtained in two charging operations are accumulated, as shown in FIG. 14B, and in the third circulation, charges are transferred to the ring-shaped CCD 521 three times; in other words, three charging operations are performed and charges obtained in the three charging operations are accumulated in the respective pixels, as shown in FIG. 14C.\nWhen the charges accumulated in the ring-shaped CCD 521 do not reach a predetermined level (level in which distance measurement can be performed on the basis of the charges), i.e., incoming light to the photoelectric conversion devices X, Y and Z is low, the number of circulations, i.e., the number of charging operation, is increased, and the charges are sequentially transferred to the ring-shaped CCD 521 and accumulated until charges are accumulated to the necessary (predetermined) level. In this manner, it is possible to obtain charges of good S/N ratio.\nWhereas, in a case where an amount of charge in the ring-shaped 521 exceeds a predetermined level within a predetermined number of circulations, i.e., in a case where incoming light to the photoelectric conversion devices X, Y and Z is high, it is necessary to adjust the amounts of charges to be stored in the pixels of the linear CCD 524 in one charging operation in order to prevent the pixels from being saturated.\nAs for adjusting the amount of charges, there are a method of adjusting a charging period by controlling an electrical shutter function, and a method for controlling a frequency for operating the photoelectric conversion devices X, Y and Z, thereby controlling a charging period.\nMore specifically, in the method of adjusting the charge amounts by controlling the electrical shutter function, if a reference charging period is 100%, then the charging period is reduced to 70%, 50%, and so on, when the object 515 is bright.\nFurther, in the method of adjusting the charge amount by controlling the frequency for operating the photoelectric conversion devices X, Y and Z, if any of the ON-pixels 522x, 522y, and 522z and the OFF-pixels 523x, 523y, and 523z is saturated when the photoelectric conversion devices X, Y and Z are operated at 1 MHz, then by operating the photoelectric conversion devices X, Y and Z in the doubled frequency, namely at 2 MHz , it is possible to halve the duration of the charging period without changing other charging conditions.\nFurther, means for changing the frequency for operating the linear CCD 524 and the ring-shaped CCD 521 depending upon the amount of external light in order to prevent the linear CCD 524 from being saturated is disclosed in the Japanese Patent Application Laid-Open No. 9-229676.\nBy adjusting the amount of charge as described above, the pixels of the linear CCD 524 are prevented from being saturated.\nIn such a hybrid-type distance measuring apparatus, a precise result of distance measurement is obtained since defects of the respective modes (active mode and passive mode) in the distance measuring operation are compensated by each other. The effects of a hybrid-type distance measuring apparatus are disclosed in the Japanese Patent Application Laid-Open No. 9-229674 by the present applicant, for instance. However, in this apparatus, since the distance measuring operations are performed both in the active mode and in the passive mode in time serially; therefore, it takes a longer time to perform the distance measuring operation comparing to a distance measuring operation performed in only one mode. Accordingly, a plurality of techniques for reducing time required for performing distance measuring operation in a hybrid-type distance measuring apparatus have been proposed. One of those techniques is disclosed in the Japanese Patent Application Laid-Open No. 9-229681 by the present applicant. According to the technique, in a case of performing a plurality of distance measuring operations consecutively, such as a case of performing a multi-point distance measuring operation, the frequency for operating a sensor and an open period of an electronic shutter set in the first operation are used in the subsequent distance measuring operations.\nIn a distance measuring apparatus mounted on a camera having a wide angle lens, for instance, a plurality of sensor arrays are arranged in the direction of the base line, as shown in FIG. 15, so as to measure distances to objects in a wide range by the respective sensor arrays.\nFIG. 15 is a view for explaining the conventional technique as described above. In FIG. 15, reference numeral 50 denotes a light-receiving lens; and 51 to 53, sensor arrays for receiving light. In FIG. 15, a light-receiving unit of a distance measuring apparatus is mainly shown, and a light-emitting device is not shown. The sensor arrays 51 to 53 receive light incoming from the right, center, and left, respectively, via the light-receiving lens 50. Reference numeral 54 denotes a selector for selecting one of the outputs from the sensor arrays 51 to 53, which is composed of a switch; and 55, a signal processing unit for performing distance measuring calculation using the output of the sensor array selected by the selector 54.\nAs shown in FIG. 15, since the respective sensor arrays 51 to 53 receive light from different directions, the intensity of light received in the respective sensor array 51 to 53 may differ from each other. Thus, if the operational frequency of a sensor and an open period of an electronic shutter which are set for measuring a distance to an object 222 using the central sensor array 52, are used as charging conditions for measuring distance to an object 221 using the left sensor array 51, since light is very strong it may cause the sensor array 51 saturation. Accordingly, in the multi-point distance measuring apparatus which measures distances to objects in the multiple directions as described above, it is necessary to independently set the operational frequency of the sensor and the open period of the electronic shutter for measuring distances in the different directions, which requires a considerably long time for performing distance measuring operation.\nNext, referring to FIG. 16, processes performed in integration clear gate (ICG) mode for determining charging conditions which do not cause saturation in the pixels of the linear CCD 524 by external light during a charge period are explained. Note, the ICG mode is performed by a not-shown controller, such as a CPU, while controlling the light-receiving unit 500.\nFirst, when the ICG mode is initiated, the operation frequency of the linear CCD 524 and the ring-shaped CCD 521, fc, is set to the lowest value, namely, 500 kHz in step S702.\nNext, a charging period is set to an initial value, i.e., the longest period. Since the charging period and other charging conditions in the linear CCD 524 can be changed by communicating with the controller (not shown) as described above, the controller generates communication data indicative of the initial value which makes the charging period longest, and sends the data to the light-receiving unit 500, thereby the charging period is set to the initial value in step S703.\nThereafter, the linear CCD 524 and the ring-shaped CCD 521 are cleared in step S704, then charging processing is initiated in step S705.\nAt the same time, the controller monitors a signal SKOS outputted from the light-receiving unit 500. The signal SKOS has a characteristic of changing its signal level (high and low) when the level of charges integrated in the ring-shaped CCD 521 exceeds a predetermined level during integrating charges in the ICG mode. The predetermined level is set on the basis of a value obtained by multiplying the capacitance of each pixel of the linear CCD 524 by a desired number. Therefore, by monitoring time since the charging processing is initiated until the level of the signal SKOS changes, it is possible to determine whether or not the current charging conditions are proper.\nNext in step S706, whether or not the level of the signal SKOS has changed is checked. If the level is not changed, then the process proceeds to step S707 where whether or not the number of circulations exceeds the maximum number of circulations, which is set in advance, is determined. If it does, then the integration of charge is terminated at that point, and the process is completed. Whereas it is does not, then the process returns to step S706 where the processes of the subsequent steps are performed.\nWhereas, if it is determined in step S706 that the level of the signal SKOS has changed, then whether or not the number of the circulations is equal to or less than a predetermined number of circulations (in FIG. 16, four), in other words, whether or not time elapsed before the level of the signal SKOS has changed is within a predetermined period, is determined in step S709. If the number of circulations exceeds the predetermined number, then the process is completed.\nNote that the number of circulations is counted by counting a clock IRCLK which will be explained later.\nWhereas, if it is determined in step S709 that the number of circulations is equal to or less than the predetermined number, namely, equal to or less than four, then the process proceeds to step S710 where whether the set operation frequency fc is the maximum frequency or not is determined.\nIf it is determined that the operation frequency fc is the maximum, then in step S711, whether or not the charging period which is currently set is the shortest is determined. If it is determined in step S711 that the charging period is the shortest, then the process is completed. Whereas, if it is not, the charging period is changed to a shorter period in step S712, then the process returns to step S703, and the subsequent steps are repeated. By changing the charging period to a shorter period, the changing condition that changed so that the level of the signal SKOS does not change within the predetermined number of circulations.\nFurther, if it is determined in step S710 that the operation frequency fc is not the maximum, the operation frequency fc is changed to a higher frequency in step S713, then the process returns to step S703 and the subsequent processes are repeated.\nAfter the ICG mode as described above is completed, the charging (integration) mode for integrating charges in the ring-shaped CCD 521 is performed, then the reading mode for outputting the charges integrated in the ring-shaped CCD 521 to a microcomputer, for example, is performed.\nThen, the microcomputer performs operations for calculating the distance to the object 515 on the basis of the outputted charges.\nFIG. 17 shows waveforms of signals when the operation frequency fc is changed from the lowest to the highest, further, the charging period is shortened by a predetermined period, in the ICG mode.\nIn FIG. 17, the signal IRCLK is a reference clock for changing the charging conditions, and the IRED 514 turns on and off in synchronization with the signal IRCLK. Further, the signal IRCLK is also used for counting the number of circulations.\nFurther, the signal SKOS is outputted from the light-receiving unit 500, and, as described above, the level of the signal SKOS is changed in accordance with the set charging conditions.\nWhen the ICG mode is initiated, 500 kHz is set as an initial frequency of the operation frequency fc. With 500 kHz, if the level of the signal SKOS has changed in a time equal to or less than four circulations (i.e., if the level of the signal SKOS has changed before the count of the signal IRCLK reaches five), the operation frequency fc is changed to a higher frequency (in this case, 1 MHz ), then the ICG mode is repeated.\nNext, with the operation frequency fc of 1 MHz, if the level of the signal SKOS has changed before the count of the signal IRCLK reaches five, the operation frequency fc (1 MHz ) is changed to a higher frequency (2 MHz ), then the ICG mode is repeated.\nThereafter, with the operation frequency fc of 2 MHz, if the level of the signal SKOS has changed before the count of the signal IRCLK reaches five, since it is not possible to change the operation frequency fc to a higher frequency, because 2 MHz is the maximum frequency, the charging period is shortened, and the ICG mode is repeated. Under the above conditions, the level of the signal SKOS does not change when the count of the signal IRCLK reaches four, in the example shown in FIG. 17. Accordingly, the ICG mode is completed, and the process moves to the charging mode.\nIn the above operation, it is necessary to drive the CCDs at a predetermined frequency while accumulating charges in the linear CCD 524 and the ring-shaped CCD 521 and while executing the ICG mode for determining the charging conditions. Whereas, it is not necessary to drive the CCDs at the predetermined frequency while communicating with the controller and clearing residual charges in the CCDs, for instance, and any frequency may be used as long as the CCDs are not improperly operated.\nHowever, as described above, the operation frequency fc when communicating with the controller and clearing residual charges in the CCDs are the same as the operation frequency fc when accumulating charges in the CCDs and executing the ICG mode for determining the charging conditions.\nMore specifically, in the example as shown in FIGS. 16 and 17, in the first execution of the ICG mode, the communication and the clearing of the CCDs in step S703 and S704 are performed at the operation frequency fc of 500 kHz, thereafter, charging operation is performed at the same driving frequency fc, i.e., 500 kHz, in the steps subsequent to step S705. In the next execution of the ICG mode, the communication and the clearing of the CCDs in step S703 and S704 are performed at the operation frequency fc of 1 MHz, thereafter, charging operation is performed at the same operation frequency fc of 1 MHz in the steps subsequent to step S705. In the next execution of the ICG mode, the communication and the clearing of the CCDs in step S703 and S704 are performed at the operation frequency fc of 2 MHz, thereafter, charging operation is performed at the same operation frequency fc of 2 MHz in the steps subsequent to step S705.\nTherefore, when a driving frequency is set to a low frequency for driving a sensor to accumulate charges or to determine charging conditions, for instance, communication with a controller, such as a microprocessor, and clearing of residual charges in the sensor are conventionally performed at the same low frequency. Therefore, compared to a case where the operation frequency is set to a high frequency, it takes a longer time to communicate with the controller and clear charges in the sensor, which causes low throughput."}
{"text": "The marking of plastics by laser and the welding of plastics parts by means of laser energy are known per se. Both are brought about as a result of absorption of the laser energy in the plastics material, either directly, by interaction with the polymer, or indirectly, with a laser-sensitive agent added to the plastics material. The laser-sensitive agent may be an organic dye or a pigment which, by absorbing the laser energy, produces a visible local change in color in the plastic. It may also be a compound which, when irradiated with laser light, is converted from an invisible, colorless form into a visible form. In the case of laser welding, the heating of the plastics materials as a result of absorption of the laser energy in the joining region is so great that the material melts and the two parts weld to one another.\nThe marking of production goods is becoming increasingly important in virtually all branches of industry as part of the general rationalization measures. Thus, for instance, there is a need to apply production data, batch numbers, expiration dates, product IDs, bar codes, corporate logos, etc. Set against conventional marking techniques such as printing, embossing, stamping, and labeling, laser marking is significantly quicker, on account of its contactless operation, and is more precise and is also unproblematic to apply to nonplanar surfaces. Since the laser markings are generated beneath the surface in the material, they are durable, robust and substantially more secure with respect to removal, alteration or even counterfeiting. Contact with other media, such as in the case of containers for liquids, and closures, is likewise uncritical for this reason—on the assumption, of course, that the plastics matrix is stable. Security and durability of production IDs, and also freedom from contamination, are extremely important in the case, for instance, of packaging of pharmaceuticals, comestibles, and beverages.\nLaser marking technology has proven very suitable in particular in connection with the marking of plastics. To be able to carry out efficient marking of plastics it is necessary to generate sufficient interaction between the plastic to be marked and the laser light. Here, firstly, it must be borne in mind that the energy introduced into the plastics should not be too high, since it may destroy the plastics article or its texture. Secondly, the laser beam must not pass through the plastic without significant interaction, since in that case the plastic cannot be marked.\nIn order to increase the interaction of the laser beam with the plastic, plastics are used in which absorption agents, also referred to as absorbers, have been incorporated. These absorbers may be, for example, laser-markable polymers or else pearlescent pigments and metallic effect pigments.\nIn the case of pearlescent pigments and metallic effect pigments, the irradiation of laser light results in heating of these pigments. In the immediate environment of the pearlescent pigments and of the metallic effect pigments, there is then a thermal change in the plastic, e.g., a carbonization or a foaming of the plastic, thereby allowing marking or identification of the plastics article.\nDE 197 26 136 A1 discloses the use of laser-markable polymers in the form of microground particles having a size of 0.1 to 100 μm. A disadvantage of these laser-markable polymers is that they can melt during the processing of the plastics doped with the laser-markable polymers. It is therefore necessary for the melting ranges of the incorporated laser-markable polymer and of the plastics system used to be harmonized with one another.\nDE 198 10 952 A1 discloses the use of pearlescent pigments or metallic luster pigments as absorbers in plastics. A disadvantage of the use of pearlescent pigments or metallic luster pigments or metallic effect pigments is that, in order to achieve satisfactory contrast after laser marking, the level of concentration must either be so high that it automatically entails coloring of the plastic with pearlescence or a metallic coloring.\nAccordingly, with use of pearlescent pigments and/or metallic effect pigments, it is not possible satisfactorily to achieve high-contrast laser marking without appreciable coloring in the case of pearlescent pigments (pearl luster effect) or without appreciable metallic coloration in the case of metallic luster or metallic effect pigments. In addition, metallic effect pigments and pearlescent pigments are relatively expensive.\nMoreover, the plateletlike structure of the pearlescent pigments or of the metallic luster or metallic effect pigments has the deleterious effect that the pigments, owing to their platelet structure, undergo orientation in the laminar flow that is an inevitable corollary of the operation of injection molding of the plastics mass, and this leads to flow lines or streaks in the plastics article produced.\nTo obtain a desired contrast in the laser marking of plastics, EP 1 145 864 A1 teaches using a mixture of metal powder and/or semimetal powder and of an effect pigment or two or more effect pigments based on phyllosilicates. Here again there is visible coloration or metallic coloring of the plastic, which is unwanted for clear and transparent plastics. Moreover, the pearlescent pigments likewise deleteriously produce streaks or flow lines in the plastics articles produced.\nDE 10 2004 053 376 A1 discloses colored laser markings and laser inscriptions on plastics that are based on the welding of a polymeric inscription medium to the plastics surface. The energy absorbers suitable for marking that are mentioned in this specification include spherical metal powders, among other absorbers. No details, however, are given concerning the size of the metal powders.\nIn accordance with the teaching of DE 10 2004 045 305 A1, the problem that exists in the prior art of the absorbers sustainedly coloring the plastics to be inscribed can be eliminated by incorporating a boride compound, preferably lanthanum hexaboride, into the plastics material. A disadvantage is that these boride compounds, especially lanthanum hexaboride constitute a significant cost factor. Consequently these boride compounds are not suitable as a laser marking agent for widespread use.\nIn order to allow marking of transparent plastics materials without coloration, use is made, in accordance with the teaching of U.S. Pat. No. 6,693,657 B2 and also of WO 2005/047009, of a laser marking agent which comprises a mixture of antimony oxide and tin oxide. WO 2005/084956 describes high-transparency plastics materials which are laser-markable and/or laser-weldable by means of nanoscale indium-tin oxide or antimony tin oxide particles. A disadvantage is that, like any other antimony compound, antimony oxide is highly toxic. Consequently this laser marking agent poses, on the one hand, a considerable risk to environment and people, both in production and processing and in disposal, since, firstly, antimony or compounds containing antimony must be used and, lastly, the plastics articles which contain antimony and/or compounds containing antimony must be disposed of again.\nWO 2002/055287 A1 describes a process for producing laser-welded composite moldings. It mentions metal flakes and metal powders as fillers. These fillers, however, are used at relatively high concentrations of 1% to 60% by weight, based on the plastics molding."}
{"text": "1. Field of the Invention\nThe present invention relates to a chip-packaging structure, and more particularly to a stacked chip-packaging structure.\n2. Description of the Related Art\nWith the progress of highly developed integrated circuit (IC) technologies, integration of internal circuits of IC chips steadily increases. As a result, number of transistors held in internal circuits of IC chips increases relatively, and sectional area of wires in the internal circuits of IC chips reduces accordingly. As the packaging technologies being continuously improved, various chip-packaging structures have been developed, including chip-packaging structures of ball grid array (BGA) structure and of multi-chip packaging module. Wherein, the BGA structure has advantages of high pin count and high reliability, while the multi-chip packaging module has shorter transfer route and better electronic properties. Such technical features make it possible to further reduce the space that the chip-packaging structure occupies, and therefore these packaging technologies could be widely applied in production of electronic devices.\nFIG. 1 is a schematic view of a conventional chip-packaging structure. The chip-packaging structure 100, through wire bonding for example, consists essentially of a substrate 110, a chip 120, a plurality of wires 130, a molding compound 140, and a plurality of solder balls 150. The substrate 110 is made of, for example, ceramic, glass, or plastics. The substrate 110 contains alternately stacked multiple conducting layers (not shown) and multiple insulating layers (not shown), while two adjacent conducting layers are separated by an insulating layer, and the conducting layers are electrically connected one another through conductive via or plating through hole. In addition, the top surface 110a and the bottom surface 110b of the substrate 110 have respectively a plurality of upper contacts 112 and a plurality of lower contacts 114, while the upper contact 112 is electrically connected to the lower contact 114 through the conducting layers in the substrate 110. Moreover, a chip 120 is disposed on the top surface 110a of the substrate 110, while the chip 120 has an active surface 122, and a plurality of bonding pads 124 (usually aluminum pads) are disposed on the periphery of the active surface 122. The bonding pads 124 are connected, through wire bonding by the wires 130, to the contacts 112 on the substrate 110, wherein the wires 130 are made of gold, for example.\nThe wires 130, the chip 120, and the upper contacts 112 on the substrate 110 are enclosed by the molding compound 140 for the protection of the chip 120 and the wires 130. The solder balls 150, made of Sn/Pb alloy for example, are disposed on the lower contacts 114. Wherein, the solder balls 150 serve as contact points of the chip-packaging structure 110 to contact with external electronic devices (e.g., printed circuits or motherboards). The solder balls 150 are disposed in a surface array distribution on the bottom surface 110b of the substrate 110, so that the chip-packaging structure 100 becomes a BGA packaging structure with high pin count.\nFIG. 2 shows a conventional multi-chip packaging module. In the packaging structure, a plurality of chip-packaging units, 200(a) and (b) (only two units are shown), are disposed on the same carrier 210, wherein the chip-packaging units 200 are of a BGA type for example, and are sequentially disposed on the carrier 210. Therefore, the chip-packaging units 200(a) and (b) can be electrically connected by internal wires (not shown) of the carrier 210, so as to construct a packaging structure of multi-chip module, such as a dynamic random access memory (DRAM) module.\nIt is worthy of notice that, when the number of the chip-packaging units increases, the requisite supporting area on the carrier increase accordingly, which makes, for example, the assembled multi-chip packaging module to have a long bar or a flat plate structure. Therefore, it is a challenge to develop appropriate chip-packaging technologies to create certain chip-packaging structures containing more chip-packaging to meet the needs for fabricating a chip-packaging structure with high pin count and a structure of multi-chip packaging module."}
{"text": "U.S. application Ser. No. 08/960,149 filed on Oct. 29, 1997, the disclosure of which is herein incorporated by reference, describes a method of generating molecular iodine in the stomach of a human or animal using an iodine reductant and an iodine oxidant. One reductant/oxidant pair identified as one of the preferred embodiments in this invention is iodide and iodate. One reason that iodide and iodate are preferred is that both of these compounds are used separately as food additives and both have GRAS (generally regarded as safe) status in the U.S. The literature describes the stability of iodide and iodate individually but there are no published reports of the stability of these two compounds in combination. In fact, these two chemicals in an appropriate medium can react rapidly with each other to yield molecular iodine.\nThis invention describes methods and compositions to stabilize an oral pharmaceutical composition that generates molecular iodine in situ wherein the active agents in said compositions are iodide and iodate. Methods for preparation of said stabilized pharmaceutical compositions are also described. It is a further objective of this invention to set forth the conditions and materials necessary to formulate a stable dosage form of iodide and iodate in the presence of accepted pharmaceutical excipients. It is a further objective of this invention to identify methods and compositions that comply with the strict stability objectives embodied in published regulatory guidelines of the FDA.\nRegulatory bodies in different territories around the world have published guidelines that define requirements that must be met to manufacture and distribute a new pharmaceutical agent for human and animal use. In the U.S. the Food and Drug Administration (FDA) regulates drugs. Stability testing is a mandated aspect of the FDA drug development process and is required for manufacture and distribution of drugs."}
{"text": "This invention relates to the field of network management, and in particular to a method and system for analyzing and verifying compliance with security policies among devices in a network.\nNetwork security is a continuing concern among network users and network managers. In a secure computer network, to prevent unauthorized access to the network, only a specified set of systems, services, and applications are permitted to use the network. Additionally, within a network, not all systems, services, and applications are permitted to communicate with each other. The term ‘security policy’ is used to identify the connections that should be permitted or denied between network elements.\nCertain devices within a network, such as routers and firewalls, can be configured to enforce security policies within the network, and complex security policies may require that multiple devices be compatibly configured.\nVerifying that each device in a network is properly configured to effect the intended security policies is a daunting task, particularly in complex networks that include multiple communication paths among network elements. This task is further complicated by the fact that different devices may use different schemes for effecting the same security policy. A security analyst must typically be aware of a variety of vendor-dependent configuration formats and protocols, and must be aware of each configuration setting's effect on the security policy of each connection to the device.\nOften, security policy analysis and verification is a ‘hit or miss’ process, with minimal analytical basis. A security analyst reviews configuration settings, corrects any obvious security loopholes, and then verifies security by attempting to send messages to or from a variety of elements on the network. If a message traverses a denied connection, or is blocked from a permitted connection, the cause of the error is determined, the configuration is modified, and the test is repeated.\nIt is an objective of this invention to provide a method and system for identifying and verifying compliance with security policies within a network. It is a further objective of this invention to provide a method and system for identifying network configuration data that is inconsistent with intended security policies.\nThese objectives, and others, are achieved by a method and system that uses a security policy database that identifies the intended security policies within a network to generate test traffic that is configured to test each defined security policy, and simulates the propagation of this traffic on a model of the network. The model of the network includes the configuration data associated with each device, and thus, if devices are properly configured to enforce the intended security policies, the success/failure of the simulated test traffic will conform to the intended permit/deny policy of each connection. Differences between the simulated message propagation and the intended security policies are reported to the user, and diagnostic tools are provided to facilitate identification of the device configuration data that accounts for the observed difference. Additionally, if a network's current security policy is unknown, test traffic is generated to reveal the actual policy in effect, to construct a baseline intended security policy.\nThroughout the drawings, the same reference numerals indicate similar or corresponding features or functions. The drawings are included for illustrative purposes and are not intended to limit the scope of the invention."}
{"text": "1. Field of the Disclosure\nThe present invention relates to fuel systems for internal combustion engines and more particularly to fuel systems that provide both primary fuel and alternate fuel to an internal combustion engine.\n2. Technical Background\nFor a variety of reasons, most of which relate to environmental and alternative energy pursuits, certain primary fuel engines, such as diesel engines, can be made to accommodate alternate fuels, such as straight vegetable oil (SVO), “biodiesel,” and other fuel oils (hereinafter “alternate fuels”). However, practical applications have demonstrated that fulltime use of these alternate fuels is not desirable. For instance, at low temperatures the viscosity of some of these alternate fuels is not optimal for use in an engine designed to run on primary fuel. The alternate fuel should be heated up to a temperature that will allow it to easily pass through the fuel delivery system. If left in an idle engine to cool, the alternate fuel also has a tendency to increase in viscosity and thus congeal or solidify inside the fuel delivery components. Thus, the alternate fuel should not be used until it has reached an appropriate temperature and it should be purged from the engine before shutting down. As a result, some of these systems typically run on diesel fuel during start-up and before shut-down, and rely on the consumer to manually switch to the alternate fuel source in between, i.e., when the alternate fuel is at a usable viscosity. For example, in an automobile that has been modified to include diesel and alternate fuel tanks, a driver can selectively toggle between the two fuels during operation of the vehicle using a switch that controls the flow of the two fuels.\nRecently, some aftermarket automated control systems have been designed. These systems control the switching between a diesel fuel source and an alternative fuel source based on the temperature of the alternate fuel. Typically, when the engine is cold, such as during engine start, the engine operates using primary fuel, and once a predetermined alternate fuel temperature or a certain vehicle speed is achieved, the engine is then switched to the alternate fuel. If low temperature operation reoccurs, the engine is then switched back to primary fuel, and so on.\nTo automate this fuel selection process, various conventional aftermarket control systems have been designed that acquire and monitor alternate fuel temperature by using various measurement devices, such as temperature sensors. However, these aftermarket systems do not address the condition whereby the engine is prematurely shut down during or a short time after alternate fuel operation. Various concentrations of alternate fuel will remain in fuel delivery components. If these fuel delivery components are subsequently allowed to cool, the alternate fuel will congeal or solidify. Fuel flow therefore cannot be reestablished for start-up.\nTherefore, an improved system and method of automatically purging an internal combustion engine of alternate fuel are desired."}
{"text": "Conventionally, electric toothbrushes that operate with a vibration that is appropriate for removing plaque have been proposed. It is desirable for a user of an electric toothbrush that information relating to where plaque is attached, whether plaque is removed by brushing, and the like is provided.\nIn order to fulfill this desire, in Patent Literature 1 (JP 2001-145645A) and Patent Literature 2 (JP 2001-170084A), a video scope and a toothbrush are integrally provided, and an image of teeth captured during brushing is displayed.\nAlso, in Patent Literature 3 (JP 2001-218624A), an image of teeth that are being brushed is captured with a camera incorporated in the brush portion of a toothbrush, and the image thereof is displayed.\nAlso, Patent Literature 4 (JP 2008-119154A) proposes a function of collecting attached substances removed from the oral cavity with a cleaning tool, and facilitating the checking of whether attached substances have been removed.\nAlso, in Patent Literature 5 (JP 2008-532619A), it is possible to adjust the power or operations of a toothbrush in accordance with information received by a sensor. Furthermore, Patent Literature 5 describes a configuration in which light is emitted from a toothbrush in order to carry out optical whitening and the like using a photochemical effect.\nWith these Patent Literatures, it is not possible to check an area unbrushed by a toothbrush, that is, a tooth where plaque remains. In contrast, the homepage of Non-Patent Literature 1 “Einstein Stella plaque” ((RF Co., Ltd.) http://www.rfsystemlab.com/product/dental/ein_plaque/index.html) and the homepage of Non-Patent Literature 2 “Intraoral Camera “MIHARU”” (http://www.shinwa-musen.co.jp/camera3/miharu.html) propose a wireless intraoral camera. The intraoral camera emits light having a predetermined wavelength to which only plaque and resin react, receives the light reflected therefrom, and converts signals of the received light into video signals to be output as video images. The intraoral camera thus has a function of providing information on an unbrushed area of teeth in video image form."}
{"text": "The present invention relates to a boot jack for gripping a boot to assist a user in removal of the user's foot therefrom.\nBoot jacks, of the general type having a frame and a plate, secured to the frame, having a U-shaped notch for receiving and gripping a wearer's boot for removal, are well known. In the simplest form, the plate is upwardly and forwardly inclined. The wearer places one foot on the back of the plate and inserts the boot on the other foot in the U-shaped notch. An upward movement of the leg, and a simultaneous downward pivoting of the user's foot about the ankle to align the foot with the leg, enables the foot to be removed from the boot. Such boot jacks may be used for cowboy style boots, but are also useful for other types of boots including overboots worn by persons in wet or snowy conditions, worn directly on the wearer's feet or over shoes on the feet.\nPatents of general background interest describing and illustrating boot jacks are U.S. Pat. No. 835,147 of H. Weber issued Nov. 6, 1906; U.S. Pat. No. 624,473 of C. Fredman et al issued May 9, 1899; U.S. Pat. No. 2,876,942 of M. Johnson issued Mar. 10, 1959; U.S. Pat. No. 1,178,109 of J. H. Staubly issued Apr. 4, 1916; Swedish Patent No. 166,129 published Feb. 3, 1959; U.S. Pat. No. 606,947 of H. Strate issued Jul. 5, 1898; U.S. Pat. No. 3,490,661 of K. L. Williams issued Jan. 20, 1970; U.S. Pat. No. 3,526,346 of L. Drouillard et al issued Sep. 1, 1970; U.S. Pat. No. 5,050,784 of Turner issued Sep. 24, 1991; Swiss Patent No. 21010 of Ernest Meyer dated Feb. 14, 1900; German Patent No. 59346 of W. Sonnemann, dated Oct. 14, 1891 and British Patent No. 664,096 of Krohnke published Jan. 2, 1952.\nIn most of these prior art boot jack and related devices, one leg must be placed behind the other, for removal of a boot, and/or there is an upward motion of the leg, controlled by muscles at the hip and knee joints, as illustrated in FIGS. 1a, b and c, to permit removal of the boot. In addition, in many previous bootjacks, the surface the foot undergoing boot removal is always flat and horizontal. This restricts the range of plantar flexion to a \"tuber angle\" of approximately 33.degree.-40.degree. . More particularly, these devices only utilize muscles which act upon the hip and knee joints. While these joints are being flexed, in order to pull the foot vertically upwards, the ankle and metatarsophalangeal joints become passively plantar flexed as they are pushed against the roof of the boot, which is held in a fixed position.\nU.S. Pat. No. 5,086,959 of Jerry issued Feb. 11, 1992 describes and illustrates a boot jack of a type having a wire frame defining a V-shaped notch having a downward incline towards the front. Again a user must place one leg behind the other, for holding the jack and a similar upward leg and foot motion, as previously described herein, is required for removal of one's foot.\nBritish Patent No. 1,001,980 of A. E. Clements Limited published Aug. 18, 1965 describes and illustrates a boot jack having upwardly and inwardly oriented converging walls defining a U-shape, for gripping a user's foot, and a platform outwardly extending to the sides and behind these walls. While the user can place one foot beside the other during removal of a boot from one foot, again a similar upward motion of the leg and knee is required in removal of the foot from the boot. A somewhat similar platform construction is taught in Swiss Patent No. 1,220 of Lingner & Kraft dated Jul. 12, 1889. In addition, U.S. Pat. No. 5,086,959, British Patent No. 1,001,980, and Swiss Patent No. 1,220 all do not utilize a leverage mechanism of the foot about the ankle joint.\nIt is an object of the present invention to provide an alternative construction of boot jack which is simple to construct and yet easy and effective in operation for removal of a users foot from a boot."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to network security.\n2. Description of the Related Art\nSome network protocols, such as certain versions of Hypertext Transfer Protocol (HTTP), allow for the transparent use of compression and encoding algorithms in sessions between a client and server. Because certain data, such as web pages, tend to compress well, the use of compression may significantly improve network performance, while the use of encoding (e.g., in encryption algorithms) is important in maintaining security.\nAs a result, however, current network security devices, such as intrusion protection systems (IPSs) face the problem that a server can send malicious traffic compressed and/or encoded to a “victim” client. Writers of such malicious content can exploit the fact that many IPS devices cannot inspect encoded or compressed HTTP sessions. In order to inspect such traffic, an ISP would have to first decompress/decode the traffic for inspection, which may not be feasible at high speeds.\nAccordingly, what is needed is an improved technique for inspecting server to client traffic in a session."}
{"text": "The present invention relates to a magnetic element such as an conductor, a choke coil, a transformer, or the like in electronic equipment, particularly a miniature magnetic element used under a large current and to a method of manufacturing the same.\nWith the reduction in size and thickness of electronic equipment, the reduction in size and thickness of components and devices used therein also has been demanded strongly. On the other hand, LSIs such as a CPU and the like have come to be made up of an increasing number of circuit components and a current of several amperes to several tens of amperes may be supplied to a power circuit provided in the LSIs. Therefore, similarly an inductor such as a choke coil used therein has been required to reduce its size, to lower the resistance, although being contrary to the size reduction, by enlarging the cross-sectional area of a coil conductor, and not to lower the inductance greatly with DC bias. The operation frequency has come to be higher and therefore it has been required that the loss in a high frequency area is low. Furthermore, in order to reduce the cost, it has been necessary that component elements with simple shapes can be assembled in easy processes. In other words, it has been demanded that a miniaturized thin inductor that can be used under a large current and at a high frequency is provided at a low cost.\nIn the case where an inductor is formed by providing a winding around a toroidal core, the inductance of the inductor is expressed by the following formula:\nLxcx9cxcexcxc3x97S=N2/r\nwherein L indicates inductance, xcexc magnetic permeability, S the cross-sectional area of a magnetic path, N the number of turns, and r the length of the magnetic path. From this formula, it is understood that a large value of L is obtained when the magnetic permeability xcexc, the cross-sectional area S of a magnetic path, and the number of turns N are increased and the length r of the magnetic path is reduced. However, when the magnetic permeability is increased, the magnetic flux density is saturated even at a small current value. The magnetic permeability is decreased at higher current values, thus deteriorating the DC bias characteristics (the inductance value (L) characteristics dependent on a direct current). The enlargement of the cross-sectional area of the magnetic path is contrary to the size reduction and results in a long lead wire in the case of the same number of turns, thus causing a high resistance. The use of a lead wire with a large cross-sectional area to prevent this further goes against the size reduction. The increase in number of turns is contrary to the size reduction and also causes a high resistance. To shorten the magnetic path leads to the size reduction but the number of turns cannot be increased in that case. Therefore, generally it has been difficult to obtain a miniature inductor that has a high inductance, excellent DC bias characteristics, and a low resistance in a winding and that can be used not only at low frequencies but also at high frequencies.\nAn inductor that has been used practically will be described as follows.\nIn an EE-type or EI-type ferrite core and a coil that have been used most commonly, because a ferrite material has a relatively high magnetic permeability and a lower saturation magnetic flux density compared to that of a metallic magnetic material, the inductance is decreased greatly due to the magnetic saturation when the ferrite material is used without being modified, resulting in poor DC bias characteristics. Therefore, in order to improve the DC bias characteristics, usually such a ferrite core and a coil have been used by providing a gap in any position in a magnetic path of the core to decrease the apparent magnetic permeability.\nIn an inductor in which a Fexe2x80x94Sixe2x80x94Al based alloy, a Fexe2x80x94Ni based alloy, or the like that has a higher saturation magnetic flux density than that of ferrite is used as a core material, because such a metallic material has a low electrical resistance, the increase in high operation frequency to several hundreds of kHz to MHz as in the recent situation results in the increase in eddy current loss and thus the inductor cannot be used without being modified. Therefore, a so-called dust core has been used, which is obtained by superposing members, which have been formed to have thin bodies, via an insulating layer or which is formed using a pulverized material that is insulated.\nIt also has been proposed to combine and use a plurality of magnetic bodies. One obtained by winding a coil around a ferrite core with rib and then dipping them into a mixed solution of magnetic powder and a resin material (JP-A-61-136213) and one obtained by preparing two members formed through the superposition of a plurality of thin magnetic metal bodies, providing a planar coil between the two members, and fixing magnetic powder with a dispersed adhesive (JP-A-9-270334) have been described as being effective for reducing the size of an inductor. In addition, one obtained by providing a planar coil between two ferrite sheets and fixing ferrite powder with a dispersed adhesive in order to reduce the leakage flux has been proposed (JP-A-6-342725), although it is not described as achieving the size reduction.\nWith respect to the configurations of inductors, many conventional inductors have been formed of an EE or EI type core and a coil. However, in order to obtain a thin inductor, JP-A-9-92540 describes using one formed by winding the coil spirally in a plane. Further, JP-A-9-205023 describes that the terminal on the internal circumference side (hereinafter referred to as an xe2x80x9cinner terminalxe2x80x9d) of a spirally wound coil is lead out by providing a cutout in a core, so that the thickness corresponding to that of the lead wire is reduced.\nHowever, when a ferrite material is used and a gap is provided anywhere in a magnetic path to decrease the apparent magnetic permeability, there has been a problem that a core vibrates in this gap portion when being operated with an alternating current, thus generating noise.\nWhen thin metallic magnetic bodies with a high saturation magnetic flux density are superposed via insulating layers, the thin bodies that can be used at high frequencies should be formed to be sufficiently thin. Therefore, the cost increases and no complicated shape can be formed, which have been problems. Further, in order to obtain a dust core with characteristics good enough, it is necessary to make the dust core dense by the application of a very high pressure of about 10t/cm2 in a molding process. Therefore, there have been problems that a special high-strength mold is required and complicated shapes are formed with difficulty.\nIn the types disclosed in JP-A-61-136213 and JP-A-6-342725 that are included in the types in which a plurality of magnetic bodies are combined and used, a member obtained by dispersing ferrite in a resin is used. However, since there is a limitation in the filling rate of the ferrite, there has been a problem that the saturation magnetic flux density of this member is low and therefore the DC bias characteristics are poor. Furthermore, in the type disclosed in JP-A-9-270334, the kind of the magnetic body to be mixed with resin is not described, but it is necessary to prepare a member formed by superposing a plurality of thin magnetic metal bodies in all cases, resulting in a high cost. In addition, since the upper and lower surfaces of an element are formed of metallic magnetic bodies, the electrical resistance is low and therefore insulation is required, and complicated shapes cannot be formed, which also have been problems.\nThe present invention seeks to provide a magnetic element, such as an inductor, a choke coil, a transformer, or the like, that is suitable for the use under a large current in various types of electronic equipment.\nA magnetic element of the present invention includes: a composite magnetic member A containing a metallic magnetic powder in an amount of 50-70 vol. % and a thermosetting resin in an amount of 50-30 vol. %; a magnetic member B that is a ferrite sintered body or a pressed-powder magnetic body of the metallic magnetic powder; and a coil. A magnetic path determined by the arrangement of the coil passes the magnetic member A and the magnetic member B in series. The coil is embedded in the magnetic member A.\nIn the magnetic element of the present invention, it is preferable that the gaps in the coil are filled with the magnetic member A. Further, it is preferable that the coil is wound around the magnetic member B.\nIt is preferable that the magnetic member B is positioned outside the magnetic member A in which the coil is embedded. In this case, it is further preferable that a plurality of plate-like magnetic members Bare included and are spaced from one another at 500-xcexcm intervals or less, particularly at 300-xcexcm intervals or less, the magnetic member A in which the coil is embedded is arranged in the intervals, and the coil is formed of a conductor wound in a planar shape.\nFurthermore, it is preferable that an oxide insulating layer is formed on the surface of the metallic magnetic powder contained in the magnetic member A. In this case, it is further preferable that the metallic magnetic powder contained in the magnetic member A contains Fe as the main component and Al, and the oxide insulating layer on the surface of the metallic magnetic powder is an insulating layer containing aluminum oxide as the main component, which is formed by a heat treatment in the presence of oxygen.\nIn this specification, the main component denotes a constituent accounting for at least 50 wt. %.\nThe present invention also provides a method of manufacturing the above-mentioned magnetic element. A first method of manufacturing the magnetic element according to the present invention includes: preparing a paste containing magnetic powder and thermosetting resin; filling gaps around the coil with the paste; and forming the magnetic member A from the paste by curing the thermosetting resin through a treatment with heat.\nA second method of manufacturing a magnetic element according to the present invention includes: preparing a slurry containing magnetic powder and thermosetting resin; forming an uncured composite sheet from the slurry; and forming the magnetic member A from the uncured composite sheet by curing the thermosetting resin through a treatment with heat and pressure."}
{"text": "When treating certain medical conditions, it is sometimes desirable to expand a frame or other radially expandable member in an orifice or conduit of a patient's body. For example, expandable tubes called stents are commonly inserted into a natural conduit of a patient's body and expanded inside the conduit to hold the conduit in an open position. Such expandable stents can be used to expand, widen, or otherwise provide structural support to various conduits of the human body, including, for example, arteries, veins, bile ducts, the esophagus, and the colon. In other treatment procedures, prosthetic heart valves that include a frame member are implanted into the body at a treatment site (e.g., a heart valve annulus). These prosthetic heart valves can be positioned in the heart valve annulus by expanding the frame member to roughly the size of the valve annulus.\nSuch frames or stents can be self-expanding or expanded using an expansion balloon. One conventional method involves positioning a frame on a balloon of a balloon catheter, maneuvering the balloon and frame to the treatment site, and inflating the balloon with a fluid to expand the frame or stent to the desired size. Such an approach, however, can have drawbacks. For example, during the expansion of the balloon the orifice or conduit is usually at least partially, if not completely, occluded, which can cause certain undesirable effects. Accordingly, it is desirable to provide methods and delivery systems that eliminate or reduce these and other potential drawbacks."}
{"text": "Many companies, municipalities, and local governments must employ a significant amount of workers to monitor and maintain equipment that may operate at distant locations relative to the operations center of a company, municipality, or local government. For example, utility companies such as electric companies must employ workers to go out and to collect data from electric meters so that an electric company can accurately measure consumption by its customers for billing purposes. Often, electric company workers go out in to the field and must physically read analog or digital meters to collect data on electricity consumption by a customer.\nIn other cases, utility meters, such as electric meters, can be provided with low power radios that can be interrogated with hand held or vehicle-mounted reading units. In such cases, workers with these hand held or vehicle-mounted reading units must be in close proximity relative to the low power radios in order to interrogate the low power radios and to collect that data from a meter. While the low power radios on the meters increase the speed and accuracy in which data can be collected by a worker, this conventional solution still requires a worker to position himself or herself in the field in close physical proximity to the meters, usually within the range of fifty feet to any particular meter. Often, because the low power radios are required by the Federal Communications Commission (FCC) to operate with such low power and because the meters are often obstructed by buildings and other physical objects, it is necessary for the worker to establish a line-of-sight coupling with the low power radio to interrogate and read it.\nAnother problem faced by utility companies, such as electric companies, is the monitoring and control of residential equipment such as air conditioners. According to some conventional programs, customers of electric companies may voluntarily relinquish control of their air conditioning unit so that the electric company can turn a particular customer's air conditioning unit on or off depending upon peak power loads monitored by the electric company. Many electric companies use existing power lines as the communications medium for controlling the operation of such air conditioning units. There are problems associated with using existing power lines as the communications medium for controlling air conditioning units. On such problem is the amount of hardware and its associated costs to support this type of communications medium that uses power lines. Also, reliability of power line carriers (PLCs) is usually low because PLCs are more susceptible to noisy, electromagnetic environments compared to other communication mediums, such as over-the-air radio frequency communications.\nIn addition to the problems faced by utility companies with controlling residential equipment from a distance and using workers to acquire data from meters, there are other problems associated with other types of equipment that may be remotely located relative to the organization that controls and maintains the equipment. For example, municipalities usually need to employ workers that monitor and maintain parking meters. A worker is needed to collect any money received by a meter and to verify that a parking meter is working properly. If information is collected by the parking meter, such as number of vehicles parked in space per day/hour or if a vehicle is present in proximity to the meter, then such information would also need to be collected by the human worker during his or her survey of the meters.\nFor other equipment, such as traffic control devices, municipalities often employ workers to perform routine checks on equipment for malfunctions and for increasing operation efficiency of the equipment. As one example, municipalities and local governments employ staff to maintain and monitor traffic lights. Usually, such staff must observe operations of traffic lights first hand in order to optimize performance and detect any malfunctions of the traffic lights. Also, staff usually must observe traffic patterns first hand in order to set the timing of the traffic lights.\nAs another example, municipalities and/or companies also employ staff to monitor and maintain automated barriers and gates to railroad crossings and draw bridges. Usually, automated barriers and gates do not have a way to communicate their operation and any environmental conditions, such as weather and traffic flow, back to a central location.\nAccordingly, there is a need in the art for a method and system that can collect information from and provide control to remote field devices relative to a central location. There is a further need in the art for a method and system that can establish communications with a remote field device without using significant communications hardware, such as wires, cables, and/or new radio equipment."}
{"text": "Gonadotropin Releasing Hormone is a decameric peptide released from the hypothalamus. In the anterior pituitary gland, GnRH activates the GnRH receptor. Activation of the GnRH receptor triggers the release of follicle stimulating hormone (FSH) and luteinizing hormone (LH). FSH and LH stimulate the biosynthesis and release of sex steroids in the gonads of both genders.\nTypically, this stimulation and release is desirable, but certain GnRH-related disorders exist where it would be beneficial to prevent activation of the GnRH receptor. For example, inhibition of the GnRH receptor can lead to a large drop in sex steroid production, which in turn can alleviate sex hormone-related conditions such as prostate cancer, endometriosis, uterine fibroids, uterine cancer, breast cancer, ovarian cancer, testicular cancer, or primary hirsutism. Moreover, there are other situations where it would be beneficial to prevent activation of the GnRH receptor, such as during some points of the in vitro fertilization process, for example, to prevent LH surge.\nCurrently marketed GnRH therapeutics are peptides that exhibit receptor antagonism in one of two ways. The first is through GnRH receptor superagonism. The GnRH receptor, when stimulated in bursts, causes normal release of the gonadotropins, FSH and LH. Under constant stimulation, the receptor becomes desensitized and the overall effect is GnRH receptor inhibition. The superagonism process is undesirable, because inhibition via this process can take up to two weeks to have an effect in human patients. During this delay, there is often an increase in disease symptoms due to the initial hormone stimulation phase. This phenomenon is referred to as flare.\nThe second method for receptor inhibition is through direct antagonism of the GnRH receptor with peptide antagonists. This causes an immediate drop in plasma LH levels. However, as mentioned above, current pharmaceuticals that cause blockade of the GnRH receptor are all peptides. As such they are not orally bioavailable and must be administered via parenteral means such as intravenous, subcutaneous, or intramuscular injection. Thus, an orally effective GnRH antagonist would be of great benefit.\nThe present invention is directed to these and other important ends."}
{"text": "In the course of an investigation on the synthesis of optically active eremophilane sesquiterpenoids from sabinene, A. van der Gen et al. [Rec. Trav. Chim. Pays-Bas, 90, 1034 (1971)]synthesized a compound to which they attributed the structure of trans-10,10-dimethyl-tricyclo[7.1.1.0.sup.2,7 ]undec-2-en-4-one. No indication as to the properties of this ketone was however given by the cited authors, nor any suggestion was formulated by them concerning its possible use as perfume or flavour ingredient. We have now found that cis-10,10-dimethyl-tricyclo[7.1.1.0.sup.2,7 ]undec-2-en-4-one possesses useful organoleptic properties and consequently it can advantageously be used in the perfume and flavour industry."}
{"text": "1. Field of the Invention\nThe present invention relates to a system and method for treating cells at a site in the body, such as at a lens capsule of an eye, thermally and, if desired, chemically. More particularly, the present invention relates to a system and method for treating cells at a site in the body, such as at a lens capsule of an eye, by applying energy to the cells to heat the cells to a temperature which will kill the cells or impede cell multiplication without causing protein denaturation to occur in the cells, and, if desired, by further exposing the cells to a material which alters a physical characteristic of the cells to kill the cells or further impede cell multiplication.\n2. Description of the Related Art\nSeveral techniques currently exist for treating cells at a selected site in the body with heat or chemicals to kill or impede multiplication of those cells to prevent undesired cell proliferation. For example, numerous types of chemotherapy drugs exists which, when injected into a tumor or delivered systemically to a patient, attack and kill cancerous cells to prevent them from further multiplying.\nRadiation techniques can also be used to kill cancerous or other undesired cells. That is, when cells are heated to a temperature of about 5xc2x0 C. or more above the normal body temperature of 37xc2x0 C., cell death begins to occur. Applying radiation to a localized site in the body, such as a tumor or other area containing undesired cells, can heat the cells at the site to temperatures in excess of 60xc2x0 C. Such high temperatures causes a phenomenon known as protein denaturation to occur in the cells, which results in immediate cell death. Accordingly, radiation therapy has been suitable in successfully treating certain types of cancers and other diseases involving uncontrolled cell growth.\nOther types of heating techniques, such as the use of probes or catheters to provide localized heat to a site of interest also exist. Like radiation therapy, these techniques also heat the cells to a high enough temperature to cause protein denaturation in the cells to thus kill the cells quickly.\nIn addition, it is also known to use photosensitive chemicals to kill cells at certain sites of interest in the body. For example, a photosensitive chemical can be injected directly into a site of interest to expose cells at that site to the chemical. A light emitting source which emits light at a wavelength that will activate the photosensitive chemical is then focused on the site of interest. Accordingly, the light activates the photosensitive chemical that has been absorbed by or is otherwise present in the cells of interest. The activated chemical kills the cells, which thus prevents undesired cell proliferation.\nAlthough the techniques mentioned above can be suitable for preventing certain types of cell proliferation and certain sites in the body, several drawbacks with these techniques exist. For example, often the use of chemotherapy drugs alone to treat a tumor or cancerous site is insufficient to kill the undesired cells. Moreover, the chemotherapy drugs also kill many normal healthy cells along with the cancerous cells, which can adversely affect the patient\"\"s health.\nThe use of radiation in conjunction with chemotherapy can have a more detrimental effect on the cancerous cells. However, as with chemotherapy, radiation often kills normal healthy cells, such as those in front of or behind the site of interest, along with the cancerous cells. Moreover, the intense heating of the cells can cause the cells to coagulate and thus block the capillaries at the site of interest. The blocked capillaries therefore prevent chemotherapy drugs from reaching the site of interest.\nIn addition, is it not known to use the above techniques to prevent unwanted cell proliferation at certain locations in the eye, such as at the retina or at the lens capsule. For example, because the retina is very sensitive, known radiation techniques can be too severe to treat cancerous cells on, in or under the retina.\nAlso, after cataract surgery, a phenomenon known as capsular opacification and, in particular, posterior capsular opacification can occur in which the epithelial cells on the lens capsule of the eye experience proliferated growth. This growth can result in the cells covering all or a substantial portion of the front and rear surfaces of the lens capsule, which can cause the lens capsule to become cloudy and thus adversely affect the patient\"\"s vision. These cells can be removed by known techniques, such as by scraping away the epithelial cells. However, it is often difficult to remove all of the unwanted cells. Hence, after time, the unwanted cells typically will grow back, thus requiring further surgery.\nAccordingly, a need exists for a system and method for preventing unwanted cell proliferation at sites in the body, especially at sites in the eye such as the retina and lens capsule, which does not suffer from the drawbacks associated with the known techniques discussed above.\nAn object of the present invention is to provide a system and method for preventing the proliferation of unwanted cells at various sites in the body, especially at sites in the eye such as the retina and lens capsule.\nAnother object of the invention is to provide a system and method for thermally and chemically treating cells at sites in the body to kill the cells or impede cell multiplication without causing protein denaturation to occur in the cells.\nA further object of the present invention is to provide a system and method which uses an energy emitting device to heat cells at a site of interest in the body to a temperature which kills the cells without causing protein denaturation in the cells, and which also chemically treats the cells, if desired, to change a physical characteristic of the cells to thus cause cell death or impede cell multiplication.\nThese and other objects of the invention are substantially achieved by providing a system and method for treating cells of a site in the body, such as at a lens capsule of an eye. The system and method employs an energy emitting device, and a positioning device, adapted to position the energy emitting device at a position in relation to the cells at the site in the body, such as the cells of the lens capsule, such that energy emitted from the energy emitting device heats the cells to a temperature which is above body temperature and below a temperature at which protein denaturation occurs in the cells, to kill the cells or impede multiplication of the cells. The energy emitting device can include a container containing a heated fluid which heats the cells to the desired temperature. Alternatively, the energy emitting device can include a heated probe, or a probe that emits radiation, such as infrared or ultraviolet radiation, laser light, microwave energy or ultrasonic energy.\nThe system and method can further employ a material delivery device that can be unitary with or separate from the energy emitting device, an can provide a material, such as a photosensitive material, to the cells at the site of interest. A light emitting device can be controlled to direct light onto the site of interest to activate the material present at the cells to alter a physical characteristic of the cells."}
{"text": "Many different types of materials are employed as functional fluids, and functional fluids are utilized in a wide variety of applications. Thus, such fluids have been utilized as electronic coolants, diffusion pump fluids, lubricants, damping fluid, power transmission and hydraulic fluids, heat transfer fluids and heat pump fluids. A particularly important application of such functional fluids has been their utilization as hydraulic fluids and lubricants in aircraft, requiring successful operation of such fluids over a wide temperature range, and fire resistant fluids.\nFunctional and hydraulic fluids employed in many industrial applications and particularly hydraulic fluids for aircraft must meet a number of important requirements. Thus, such hydraulic fluids particularly for aircraft use, should be operable over a wide temperature range, should have good stability at relatively high temperatures and preferably have lubricating characteristics. In addition to having the usual combination of properties making it a good lubricant or hydraulic fluid, such fluid should also have relatively low viscosity at extremely low temperatures and an adequately high viscosity at relatively high temperatures, and must have adequate stability at the high operating temperatures of use. Further, it is of importance that such fluids be compatible with and not adversely affect materials including metals and non-metals such as seals of the system in which the fluid is employed. It is also important in aircraft hydraulic fluids and lubricants that such fluids have as high a fire resistance as possible to prevent ignition if such fluids are accidentally or as result of damage to the hydraulic system, sprayed onto or into contact with surfaces of materials of high temperature. Another important property for application of a hydraulic fluid in aircraft is the provision of a low density fluid to increase pay load.\nHydraulic fluids in commercial jet aircraft are exposed to temperatures ranging from below -40.degree.F to over 200.degree.F. Within these temperature extremes, it is necessary for the fluid to maintain a reasonably low viscosity when cold, and yet not become too thin when hot.\nIn presently available commercial functional or hydraulic fluids, phosphate esters are among the most commonly employed base stocks, of which tributyl phosphate and dibutyl phenyl phosphate are widely used components. Both of the latter phosphates are too thin at high temperatures, and their use alone would result in rapid wear of moving parts. Other phosphate esters, such as tricresyl phosphate, for example, which provide the requisite high temperature viscosity become too thick to be useful at low temperatures. Even mixtures of various phosphate esters such as those noted above do not provide the required viscosity characteristics at both low and high temperatures. Accordingly, it has been the practice to achieve the required wide viscosity range required for aircraft hydraulic fluids by adding to a thin base stock, such as phosphate ester or mixtures thereof, a small proportion, e.g., up to 10%, of a polymeric material such as polyalkyl acrylates or methacrylates, whose solubility characteristics in the base stock are chosen so that the polymeric material thickens the fluid more at high temperatures than at low temperatures, and thus functions as a viscosity index (VI) improver.\nHowever, during use, fluids of the above type containing polymers such as the polyalkyl acrylates or methacrylates, tend to decompose due to the shearing forces of the mechanical components encountered in an aircraft hydraulic system, on the fluid, and producing acids. This results in a high degree of acid build-up during use, which is detrimental in causing corrosion of metal surfaces with which the fluid is in contact and also causes further decomposition of the fluid.\nIn U.S. Pat. No. 3,637,507, improved acid stability of such functional or hydraulic fluids is achieved by the addition of monoepoxides and particularly certain diepoxides to hydraulic fluids containing a phosphate ester and particularly polyalkyl methacrylate viscosity index improver. However, in currently manufactured fluids of this type, the functionality of the epoxide is severly limited by the reaction which epoxides can have with acrylate and methacrylate viscosity index improver. Thus, although diepoxides such as the 3,4-epoxycyclohexylmethyl 3,4-epoxycyclohexane carboxylate disclosed in the above patent are preferred over the monoepoxide in the above fluids, the concentration of the latter diepoxide must be kept low since the reaction of polyfunctional epoxides with acrylate and methacrylate type viscosity index improvers leads to the formation of harmful insoluble deposits in aircraft hydraulic systems. Thus, it is particularly noteworthy that in all of the examples of the above patent, and employing polyalkyl methacrylate viscosity index improver, the concentration of the epoxide component, including the above-noted preferred diepoxide of the patent, is employed in a concentration ranging only from about 0.5 to about 1%. The necessity for using such low concentrations of such diepoxides in order to avoid formation of undesirable deposits in the fluid, due to reaction with acrylate and methacrylate viscosity index improvers, limits the effectiveness of such epoxides in preventing acid build-up.\nThe present invention relates to functional fluid compositions having good fire resistance and desirable viscosity characteristics at both high and low temperatures, and is particularly directed to functional fluid compositions having the above-noted properties and consisting essentially of a phosphorus compound, especially a phosphate ester, a particular viscosity index improver other than a polyacrylate or polymethacrylate, and an epoxide compound of a type which permits use of a high concentration of such epoxide in conjunction with said viscosity index improver, to effectively prevent formation of harmful insoluble deposits in the fluid while effectively functioning as an acid absorber to prevent acid build-up."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical recording medium, in particular to a write once optical recording medium capable of recording information therein and reproducing the recorded information therefrom by causing the optical characteristics, such as transmittance and reflectance, of a recording material for use in the optical recording medium to change by the application of an optical beam thereto. In addition, the present invention also relates to a method of recording information in the above-mentioned optical recording medium.\n2. Discussion of Background\nThe wavelength of a laser beam used for currently employed write once optical disk systems of a write once read many type (WORM), and a compact disc-recordable (CD-R) type is in the range of 770 nm to 790 nm, and a recording medium for the write once optical disk system is designed in such a fashion that recording and reproduction can be performed in the above-mentioned wavelength range.\nSooner or later, it will be inevitable that the capacity of such a recording medium must be significantly increased in accordance with the amount of information to be handled. It is also inevitable that the wavelength of the laser beam used for the recording and reproduction will be shortened.\nIn recent years, there have been proposed many types of write once optical disks for data recording and reproduction, using a cyanine dye or a phthalocyanine dye as a recording material.\nFor example, the write once optical disks using as the recording material a cyanine dye are disclosed in Japanese Laid-Open Patent Applications 57-62093, 58-56892, 58-112790, 58-114989, 59-85791, 60-83236, 60-89842, and 61-25886. Furthermore, the write once optical disks using as the recording material a phthalocyanine dye are disclosed in Japanese Laid-Open Patent Applications 61-150243, 61-177287, 61-154888, 61-246091, 62-39286, 63-37791, and 63-39888.\nHowever, there has not yet been developed any recording material that has high light resistance and excellent preservation stability, and is capable of performing the recording and reproducing operation using an optical pick-up of a laser beam with a wavelength of 700 nm or less.\nThe currently employed CD-R is also designed in such a fashion that recording and reproduction can be performed by a laser beam with a wavelength in the range of 770 nm to 790 nm.\nIn the case of the above-mentioned CD-R system, the increase in capacity of the recording medium also becomes inevitable, and there is the trend toward shortening of the wavelength of the laser beam to be used for data recording and reproduction.\nIn CDs and CD-ROMs which are currently employed, aluminum (Al) is coated on the concave and convex portions of the substrate therefor and the wavelength-dependence of the reflectance of Al is so small that even if the wavelength of the laser beam used therefor is shortened in the future, reproduction of information recorded in such CDs and CD-ROMs is possible.\nIn contrast to this, the recording layer of the currently employed CD-R, which contains therein a dye with a maximum absorption wavelength in the range of 680 nm to 750 nm, is designed so as to exhibit a maximum reflectance when a laser beam with a wavelength of 770 nm to 790 nm is applied thereto, with the optical characteristics and constants and the thickness and structure of the recording layer being taken into consideration. Therefore, when a laser beam with a wavelength of 700 nm or less is applied to the currently employed CD-R, the reflectance exhibited by the recording layer is extremely small, so that the currently employed CD-R cannot cope with the shortening of the wavelength of the employed laser beam in the future. As a result, it will be highly possible that the information recorded in the currently employed CD-R system cannot be reproduced therefrom by such a future system.\nMany recording materials, for example, the combination of cyanine dye/metal reflection layer, phthalocyanine dye/metal reflection layer, and azo metal chelate dye/metal reflection layer, are proposed for the conventional CD-R. To be more specific, as the recording material for the CD-R, the combination of a cyanine dye/metal reflection layer is proposed in Japanese Laid-Open Patent Applications 1-159842, 2-42652, 2-13656, and 2-168446; a phthalocyanine dye as the recording material is proposed in Japanese Laid-Open Patent Applications 1-176585, 3-215466, 4-113886, 4-226390, 5-1272, 5-171052, 5-116456, 5-69860 and 5-139044; and an azo metal chelate dye as the recording material is proposed in Japanese Laid-Open Patent Applications 4-46186, 4-141489, 4-361088, and 5-279580.\nHowever, there is no write once optical recording medium that can solve the above-mentioned problems."}
{"text": "Phototropic (also known as photochromic and photothermochromic) compounds are generally known as compounds which isomerize or tautomerize reversibly to another form when exposed to actinic light of a suitable wave length. When this wavelength of irradiation is removed and either heat or light of a different wavelength of irradiation is applied, the reaction is reversed to give the original phototropic compound.\nIt has now been discovered that a specific class of phototropic compounds, when irradiated to an aci- or acidic form, are effective as photoinitiators for cationic polymerization of various monomers resulting in polymers thereof and such photo-induced reaction is eminently suitable for use in the graphic arts or as light-cured coatings.\nVinyl ethers and epoxy resins have been photopolymerized employing various catalysts and promoters therefor. For example, U.S. Pat. No. 3,708,296 issued Jan. 2, 1973 to S. I. Schlesinger and commonly assigned herewith, employs photosensitive catalyst precursors which are aryl diazonium salts of complex halogenides which release a Lewis acid upon irradiation to initiate polymerization of epoxide monomers and prepolymers.\nAnother procedure, U.S. Pat. No. 3,347,676 employs metal salts such as complex silver borates in combination with organic halogens.\nNone of the prior procedures are without their particular difficulties however. For example, epoxides and related compositions containing photosensitive catalyst precursors have a tendency to gel on standing, often necessitating the need for gelation inhibitors to prevent premature reaction, while silver compounds are becoming increasingly expensive, and are not practical alternatives for large scale commercial production. Accordingly, it is desirable to identify and discover new monomer catalyzing agents useful in the photopolymerization process which are not subject to and overcome the deficiencies now existing in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to a chemical mechanical polish (CMP) process, and more particularly, to a CMP process and a method for improving accuracy of determining a polish endpoint thereof.\n2. Description of the Related Art\nWith advance of semiconductor technology, number of transistors which can be fabricated in a chip per unit area is increased. The higher the device integrity is, the more the unplanar surface of a chip is undesirable. As a result, enhancing surface planarization of a chip becomes essentially important. A chemical mechanical polish (CMP) process is a method to provide a uniform surface planarization, and has become a vital thin film planarization technique for deep sub-micron technology.\nUse the commonly applied tungsten plug CMP process as an example. FIG. 1A is a cross sectional view of a prior art metal interconnect structure. The silicon oxide layer 111 is formed on the substrate 101 of the chip 100. The silicon oxide layer 111 comprises the opening 130. The adhesion layer 140 covers the silicon oxide layer 111. The material of the adhesion layer 140 can be, for example, titanium nitride/titanium (TiN/Ti). The conductive layer 150 is formed over the adhesion layer 140. The material of the conductive layer 150 can be, for example, tungsten. The opening 130 serves as a contact hole. The adhesion layer 140 and the conductive layer 150 filled in the opening 130 serve as a conductive plug. Since conductive lines are formed on the surface of the silicon oxide layer 111 to connect with the conductive plugs, the adhesion layer 140 and the conductive layer 150 over the surface of the silicon oxide layer 111 should be removed by a CMP process.\nIn order to avoid residues, dishing, or erosion resulting from under-polish or over-polish by the CMP process, how to accurately detect polish endpoints of the CMP process becomes an important topic in this field. For the presently used method of determining the polish endpoint, an IR light irradiates the conductive layer 150 while the to-be-polished layer, e.g., the conductive layer 150 and the adhesion layer 140 of FIG. 1A, is polished by the CMP process. The reflection of the IR light from the to-be-polished layer is continuously detected. Most of the material of the to-be-polished layer contains metal, and the silicon oxide layer 111 is non-metal. Accordingly, the reflection of the IR light from the to-be-polished layer is much stronger than that from the silicon oxide layer 111. When the reflection is weakened, it can be determined that the to-be-polished layer is removed and the surface of the silicon oxide layer 111 is exposed. At this moment, it also indicates the polish endpoint of the CMP process.\nIn another aspect, since the device integrity is increased due to advance of semiconductor technology, the high device integrity reduces line width and increases difficulty for photolithographic process. As a result, the misalignment issue is easy to happen. The issue becomes serious when a reflective layer, such as a metal layer or a polysilicon layer, is included in the structure. The reflective layer results in reflection of incident light on the film surface and diffraction. Thus, the photolithographic pattern may be incorrectly transferred. To prevent the issue described above, an anti-reflection coating (ARC) layer or a dielectric anti-reflection coating layer (DARC) layer is formed on the film layer to reduce the errors caused by the reflection during the photoresist exposure step, and to enhance device yield.\nThough the ARC layer or the DARC layer is necessary for the photolithographic process, a new issue arises. During the CMP planarization process, the layer that should be polished is the material layer 120 of FIG. 1B, rather than the original single silicon oxide layer 111. The material layer 120 comprises, for example, the silicon oxide layer 111, the first DARC layer 113, the second DARC layer 115 and the cap oxide layer 117, which are sequentially formed on the substrate 101. These film layers enhance reflection of the IR light, which becomes irregular from the original reflection. With reducing the thickness of the to-be-polished layer, the reflection of IR light is enhanced and reflection variation becomes severe. This phenomenon may make the polish endpoint detection failed. If the film layers are under-polished, the residue of the material layer 120 may be generated on the conductive layer 150, causing bridge. In contrast, if the film layers are over-polished, the material layer 120 may be eroded. The material layer 120 may become thinner, and the conductive layer 150 is also subject to the over-polish. In a high-density device, the dielectric material among the metal interconnect structure is not sufficient to serve electrical isolation so as to affect device operation."}
{"text": "As an example of a mobile communication system to which the present invention is applicable, a 3rd Generation Partnership Project Long Term Evolution (hereinafter, referred to as LTE) communication system is described in brief.\nFIG. 1 is a view schematically illustrating a network structure of an E-UMTS as an exemplary radio communication system. An Evolved Universal Mobile Telecommunications System (E-UMTS) is an advanced version of a conventional Universal Mobile Telecommunications System (UMTS) and basic standardization thereof is currently underway in the 3GPP. E-UMTS may be generally referred to as a Long Term Evolution (LTE) system. For details of the technical specifications of the UMTS and E-UMTS, reference can be made to Release 7 and Release 8 of “3rd Generation Partnership Project; Technical Specification Group Radio Access Network”.\nReferring to FIG. 1, the E-UMTS includes a User Equipment (UE), eNode Bs (eNBs), and an Access Gateway (AG) which is located at an end of the network (E-UTRAN) and connected to an external network. The eNBs may simultaneously transmit multiple data streams for a broadcast service, a multicast service, and/or a unicast service.\nOne or more cells are present per eNB. A cell is configured to use one of bandwidths of 1.44, 3, 5, 10, 15, and 20 MHz to provide a downlink or uplink transport service to several UEs. Different cells may be set to provide different bandwidths. The eNB controls data transmission and reception for a plurality of UEs. The eNB transmits downlink scheduling information with respect to downlink data to notify a corresponding UE of a time/frequency domain in which data is to be transmitted, coding, data size, and Hybrid Automatic Repeat and reQuest (HARQ)-related information. In addition, the eNB transmits uplink scheduling information with respect to uplink data to a corresponding UE to inform the UE of an available time/frequency domain, coding, data size, and HARQ-related information. An interface may be used to transmit user traffic or control traffic between eNBs. A Core Network (CN) may include the AG, a network node for user registration of the UE, and the like. The AG manages mobility of a UE on a Tracking Area (TA) basis, each TA including a plurality of cells.\nAlthough radio communication technology has been developed up to LTE based on Wideband Code Division Multiple Access (WCDMA), demands and expectations of users and providers continue to increase. In addition, since other radio access technologies continue to be developed, new advances in technology are required to secure future competitiveness. For example, decrease of cost per bit, increase of service availability, flexible use of a frequency band, simple structure, open interface, and suitable power consumption by a UE are required."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical wavelength division multiplexing network which multiplexes and transmits optical signals having a plurality of different wavelengths.\n2. Description of the Related Art\nFIG. 25 shows an example of the constitution of a conventional optical wavelength division multiplexing network. The network shown in FIG. 25 has a ring structure comprising two or three layers. The network of FIG. 25 will be explained in separate sections comprising a network (1) 11, a network (2) 12, a network (3) 13, a network (4) 14, and a network (5) 15. The network (1) 11 has a ring structure, and is provided at the highest level. The network (1) 11 comprises at least one center node 21 and two or more (three in FIG. 1) remote nodes 22, 23, and 24. The network (2) 12 is a ring network comprising a node (#4) 24, which is one of the remote nodes of the network (1) 11, and is provided below the network (1) 11. The network (3) 13 comprises a tree-shaped structure centered on a node (#41) 25, which is one of the nodes of the network (2) 12, and is provided below the network (2) 12. The network (4) 14 comprises a ring-shaped structure centered on a node (#3) 23, which is one of the remote nodes of the network (1) 11, and is provided below the network (1) 11. The network (5) 15 comprises a tree-shaped structure centered on a node 27, which is one of a plurality of nodes 23, 26, 27, and 28 of the network (4) 14, and can conceivably be provided below the network (4) 14. Optical network units (ONU) (they are also called optical service units.) 51 to 56 comprise the subscribers of each home, business office, and the like, and are provided in the network (3) 13 or the network (5) 15.\nIn FIG. 25, the nodes of network (1) 11, the network (2) 12, and the network (4) 14, are connected by optical fiber transmission paths 60, 60, . . . which comprise a plurality of optical fibers. The ONUs are connected to the nodes of the networks (3) 13 and (5) 15 by optical fiber transmission paths 70, 70, . . . which comprise a single optical fiber. Equipment for signal termination electrically processes transmission signals, which have been converted from optical signals to electrical signals, and are provided at the nodes 21 to 28.\nIn this explanation, traffic from subscribers in the network (3) 13 or the network (5) 15 is assumed to be 1.5 Mb/s. The traffic from the subscribers is multiplexed at the subscriber office (node 25 or node 27), and transmitted to the node (#4) 24 in the network (2) 12 or the remote node (#3) 23 in the network (1) 11 at a higher transmission speed of, for example, 52 Mb/s. At the node (#41) 25 and the node 27, traffic sent from other nodes in the network (2) 12 or the network (4) 14 is combined with the multiplexed traffic from the subscribers, and transmitted to the next node in the network (2) 12 or the network (4) 14 at an even high transmission speed. Transmission speed conversion and the like is also carried out at the nodes in the network (1) 11. That is, electrical processing is carried out at each node.\nIn conventional optical networks such as that shown in FIG. 25, when starting a new high-speed access service for users at a speed of, for example, approximately 150 Mb/s or 1 Gb/s, transmission apparatuses which carry out electrical processing for multiplexing the traffic must be provided at each node, since there are several users belonging to the network (3) 13 or the network (5) 15. Consequently, the initial expenditure is considerable. Moreover, depending on the region, there may be fewer users per node, leading to a drawback of expenditure efficiency."}
{"text": "While a number of recent efforts are being made to achieve a finer pattern rule in the drive for higher integration and operating speeds in LSI devices, the commonly used light exposure lithography is approaching the essential limit of resolution determined by the light source wavelength.\nAs the light source used in the lithography for resist pattern formation, g-line (436 nm) or i-line (365 nm) from a mercury lamp has been widely used. One means believed effective for further reducing the feature size is to reduce the wavelength of exposure light. For the mass production process of 64 M-bit DRAM, the exposure light source of i-line (365 nm) was replaced by a KrF excimer laser having a shorter wavelength of 248 nm. However, for the fabrication of DRAM with a degree of integration of 1 G or more requiring a finer patterning technology (processing feature size 0.13 μm or less), a shorter wavelength light source is required, and in particular, photolithography using ArF excimer laser light (193 nm) is now under investigation.\nOn the other hand, it is known in the art that the bilayer resist method is advantageous in forming a high-aspect ratio pattern on a stepped substrate. In order that a bilayer resist film be developable with a common alkaline developer, high molecular weight silicone compounds having hydrophilic groups such as hydroxyl and carboxyl groups must be used.\nAmong silicone base chemically amplified positive resist compositions, recently proposed were those compositions for KrF laser exposure comprising a base resin in the form of polyhydroxybenzylsilsesquioxane, which is a stable alkali-soluble silicone polymer, in which some phenolic hydroxyl groups are blocked with t-BOC groups, in combination with a photoacid generator (see JP-A 6-118651 and SPIE vol. 1925 (1993), p. 377). For ArF laser exposure, positive resist compositions comprising as a base a silsesquioxane of the type in which cyclohexylcarboxylic acid is substituted with an acid labile group were proposed (see JP-A 10-324748, JP-A 11-302382, and SPIE vol. 3333-07 (1998), p. 62). For F2 laser exposure, positive resist compositions comprising a silsesquioxane having hexafluoroisopropanol as a dissolvable group as a base were proposed (see JP-A 2002-55456). The above polymer bears in its backbone a polysilsesquioxane containing a ladder skeleton produced through polycondensation of a trialkoxysilane or trihalosilane.\nSilicon-containing (meth)acrylate polymers were proposed as a resist base polymer having silicon pendants on side chains (see JP-A 9-110938, J. Photopolymer Sci. and Technol., Vol. 9, No. 3 (1996), pp. 435-446).\nThe undercoat layer of the bilayer resist process is formed of a hydrocarbon compound which can be etched with oxygen gas, and must have high etching resistance since it serves as a mask when the underlying substrate is to be etched. For oxygen gas etching, the undercoat layer must be formed solely of a silicon atom-free hydrocarbon. To improve the line-width controllability of the upper layer of silicon-containing resist and to minimize the sidewall corrugation and collapse of the pattern by standing waves, the undercoat layer must have the function of an antireflective film. Specifically, the reflectivity from the undercoat layer back into the resist film must be reduced to below 1%.\nIn the case of a subbing antireflective film used in the monolayer resist process, when a high reflection substrate such as polysilicon or aluminum underlies, a material having an optimum refractive index (n value) and absorption coefficient (k value) is designed to an appropriate film thickness, whereby the reflectivity from the substrate can be reduced to below 1%, achieving a significant antireflection effect. In an example wherein the wavelength is 193 nm and a resist has a refractive index of 1.7, if the subbing antireflective film has a refractive index (real part of complex refractive index) n of 1.5, an extinction coefficient (imaginary part of complex refractive index) k of 0.5, and a thickness of 42 nm, then the reflectivity becomes below 0.5% (see FIG. 1). However, if the substrate has steps, the antireflective film largely varies its thickness at the steps. The antireflection effect of the subbing film utilizes not only the absorption of light, but also the interference effect arising from a properly designed film thickness. The first base of 40 to 45 nm having the enhanced interference effect has an accordingly enhanced antireflection effect, but the reflectivity largely varies with a change of film thickness. JP-A 10-69072 discloses a high conformity antireflective film-forming material in which the molecular weight of a base resin is increased to minimize the variation of film thickness at steps. As the molecular weight of a base resin increases, there arise problems that more pinholes generate after spin coating, filtration becomes difficult, a viscosity change with the passage of time leads to a variation of film thickness, and crystals precipitate at the nozzle tip. The conformal behavior is developed only at relatively low steps.\nIn another method using a film thickness of at least the third base (i.e., of at least 170 nm) where the variation of reflectivity due to a film thickness variation is relatively small, the variation of reflectivity due to a film thickness variation is small and the reflectivity is kept below 1.5% as long as the k value is in a range of 0.2 to 0.3 and the film thickness is at least 170 nm. However, when the etching load that the overlying resist layer has to bear is considered, the approach of thickening the antireflective film encounters a limit, the limit of film thickening being of the order of the second base of up to 100 nm.\nIn the event the underlay below the antireflective film is a transparent film like an oxide or nitride film and steps exist below that transparent film, the thickness of the transparent film varies even if the surface of the transparent film is planarized as by chemical mechanical polishing (CMP). In this event, it is possible to make the thickness of the overlying antireflective film uniform. If the thickness of a transparent film underlying the antireflective film varies, the thickness of the minimum reflective film is shifted by the thickness of the transparent film. Even if the thickness of the antireflective film is set equal to the thickness of the minimum reflective film when the underlay is a reflective film, the reflectivity can be increased due to a variation of the thickness of the transparent film.\nThe materials of which the antireflective film is made are generally divided into inorganic and organic materials. A typical inorganic material is a SiON film. This has the advantages that it can be formed by CVD of a gas mixture of silane and ammonia, and the etching load on resist is light due to a high selective ratio of etching relative to resist, but the range of application is restricted because of difficulty of peeling. Because of a nitrogen atom-containing basic film, another drawback arises that it is susceptible to footing in the case of positive resist and an undercut profile in the case of negative resist.\nThe organic material has the advantages that spin coating is possible without a need for a special equipment as needed for CVD and sputtering, peeling is possible like resist, no footing occurs, the shape is obedient, and adhesion to resist is good. Thus a number of antireflective films based on organic materials have been proposed. For example, JP-B 7-69611 describes a composition comprising a condensate of a diphenylamine derivative with a formaldehyde-modified melamine resin, an alkali-soluble resin, and a light-absorbing agent; U.S. Pat. No. 5,294,680 describes the reaction product of a maleic anhydride copolymer with a diamine light-absorbing agent; JP-A 6-118631 describes a composition comprising a resin binder and a methylol melamine type heat crosslinking agent; JP-A 6-118656 describes a composition based on an acrylic resin having a carboxylic group, an epoxy group and a light-absorbing group within a common molecule; JP-A 8-87115 describes a composition comprising a methylol melamine and a benzophenone light-absorbing agent; and JP-A 8-179509 describes a polyvinyl alcohol resin with a low molecular weight light-absorbing agent added thereto. In all these patents, a light-absorbing agent is either added to a binder polymer or introduced into a polymer as substituent groups. However, since most light-absorbing agents have aromatic groups or double bonds, the addition of a light-absorbing agent undesirably increases dry etching resistance and rather reduces a selective ratio of dry etching relative to the resist. As the feature size becomes finer, the drive toward resist film slimming is accelerated. In the ArF exposure lithography of the next generation, acrylic or alicyclic polymers are used as the resist material, indicating that the etching resistance of the resist becomes poor. A further consideration is the problem that the thickness of antireflective film must be increased as mentioned above. Then, etching is an acute problem. There is a need for an antireflective film having a high selective ratio of etching relative to resist, that is, a high etching speed.\nThe function required as an antireflective film for the undercoat layer of the bilayer resist process differs from the function required in the monolayer resist process. Since the undercoat layer of the bilayer resist process serves as a mask when the substrate is etched, it must have a high etching resistance under the substrate etching conditions. As opposed to the antireflective film in the monolayer resist process required to have a high etching speed in order to mitigate the load to the monolayer resist, the inverse performance is required. To provide sufficient substrate etching resistance, the thickness of the undercoat layer must be equal to or greater than the monolayer resist, that is, as thick as 300 nm or greater. At a film thickness of at least 300 nm, the variation of reflectivity due to a film thickness change is substantially converged so that the antireflection effect due to phase difference control is no longer expected.\nNow, the results of calculation of reflectivity at varying film thicknesses up to the maximum thickness of 500 nm are shown in FIGS. 2 and 3. Assume that the exposure wavelength is 193 nm, and the upper layer resist has an n value of 1.74 and a k value of 0.02. FIG. 2 shows substrate reflectivity when the k value of the undercoat layer is fixed at 0.3, the n value varies from 1.0 to 2.0 on the ordinate and the film thickness varies from 0 to 500 nm on the abscissa. In an imaginary case wherein the undercoat layer of the bilayer resist process has a thickness of 300 nm or greater, optimum values at which the reflectivity is reduced to or below 1% exist in the refractive index range of 1.6 to 1.9 which is approximate to or higher than that of the upper layer resist.\nFIG. 3 shows substrate reflectivity when the n value of the undercoat layer is fixed at 1.5 and the k value varies from 0.1 to 0.8. In the k value range of 0.24 to 0.15, the reflectivity can be reduced to or below 1%. By contrast, the antireflective film in the form of a thin film of about 40 nm used in the monolayer resist process has an optimum k value in the range of 0.4 to 0.5, which differs from the optimum k value of the undercoat layer in the bilayer resist process having a thickness of 300 nm or greater. For the undercoat layer in the bilayer resist process, an undercoat film having a lower k value, that is, more transparent is necessary.\nAs the material for forming an undercoat layer for 193 nm exposure, copolymers of polyhydroxystyrene with acrylic compounds are under study as described in SPIE vol. 4345, p. 50 (2001). Polyhydroxystyrene has a very strong absorption at 193 nm and its k value is as high as around 0.6 by itself. By copolymerizing with an acrylic compound having a k value of almost 0, the k value is adjusted to around 0.25.\nHowever, the resistance of an acrylic compound to substrate etching is weak as compared with polyhydroxystyrene, and a considerable proportion of the acrylic compound must be copolymerized in order to reduce the k value. As a result, the resistance to substrate etching is considerably reduced. The etching resistance is not only reflected by the etching speed, but also evidenced by the development of surface roughness after etching. Through copolymerization of acrylic compound, the surface roughness after etching is increased to a level of significance.\nNaphthalene ring is one of rings that have a more transparency at 193 nm and a higher etching resistance than benzene ring. JP-A 2002-14474 proposes an undercoat layer having a naphthalene or anthracene ring. However, since naphthol-copolycondensed novolac resin and polyvinyl naphthalene resin have k values in the range of 0.3 to 0.4, the target transparency corresponding to a k value of 0.1 to 0.3 is not reached, with a further improvement in transparency being necessary. The naphthol-copolycondensed novolac resin and polyvinyl naphthalene resin have low n values at 193 nm, as evidenced by a value of 1.4 for the naphthol-copolycondensed novolac resin and a value of only 1.2 for the polyvinyl naphthalene resin as long as the inventors measured. JP-A 2001-40293 and JP-A 2002-214777 describe acenaphthylene polymers which have lower n values and higher k values at the wavelength of 193 nm than at 248 nm, both falling outside the target values. There is a need for an undercoat layer having a high n value, a low k value, transparency and high etching resistance.\nJP-A 6-202317, JP-A 8-179502, JP-A 8-220750, JP-A 8-292565, and JP-A 9-15855 disclose i-line resists comprising a copolycondensed polymer of cresol and dicyclopentadiene as a base resin. Copolymerization with dicyclopentadiene was under study to develop higher transparency novolac resins. JP-A 10-282666 proposes a curable resin in the form of a resol-dicyclopentadiene copolycondensed polymer having pendant glycidyl groups. Further, JP-A 6-80760 and JP-A 7-5302 propose a resist composition comprising naphthol condensed with aldehyde."}
{"text": "1. Technical Field\nThe present disclosure relates to a chemical sensor formed in a semiconductor die.\n2. Description of the Related Art\nChemical sensors are used in a variety of applications. Chemical sensors are used in medical applications, industrial applications, automotive applications, security applications, and domestic applications. Some examples of chemical sensors are blood glucose sensors, carbon dioxide detectors, automobile exhaust emission monitors, radon detectors, carbon monoxide detectors, explosives detectors, and a large variety of other applications.\nIn the past, many chemical sensors have been large and relatively expensive. Some chemical sensors are used in applications in which they may only be used a single time. Such single use sensors are typically used in biomedical applications. It can be very expensive to replace an entire chemical detection system after each use. It may thus greatly increase the cost of medical diagnosis and treatment."}
{"text": "CD27, a TNF receptor family member was identified as a membrane molecule on human T cells (van Lier et al., 1987, J Immunol 139:1589-96). According to current evidence, CD27 has a single ligand, CD70, which is also a TNF family member (Goodwin et al., 1993, Cell 73:447-56). CD27 and CD70 have also been identified and cloned in the mouse system (Gravestein et al., 1993, Eur J Immunol 23:943-50; Tesselaar et al., J Immunol 159:4959-65).\nCD27 is exclusively expressed by hematopoietic cells, in particular those of the lymphocyte lineage, i.e. T-, B- and NK cells. In the human system, CD27 expression in the αβ T cell lineage is induced during positive selection of thymocytes and maintained in naive conventional CD4+ and CD8+ T cells (Vanhecke et al., 1997, J Immunol 159:5973-83). Upon T cell activation via the T cell antigen receptor (TCR), CD27 expression increases, in a transient manner (van Lier et al., 1987, J Immunol 139:1589-96). Next, CD27 is shed from the surface of activated T cells and the soluble form of CD27 can be detected in the serum marker for (chronic) T cell activation (Hintzen et al., 1991, J Immunol 147:29-35). Among peripheral T cells, permanent loss of CD27 expression results from persistent antigenic stimulation and hallmarks terminally differentiated effector/memory T cells, while central memory T cells maintain CD27 (Baars et al., 1995, J Immunol 154:17-25; Hamann et al., 1997, J Exp Med 186:1407-18). CD27 is also expressed on human γδ T cells and induced during thymic development (Offner F et al, 1997, J Immunol 158:4634-41). Moreover, loss of CD27 expression is a hallmark of chronically stimulated γδ T cells (Gioia C et al., 2002, J Immunol 168:1484-9). Generally, CD27 is an exquisite marker for cellular activation and differentiation stages and used as such in human clinical diagnostics and research.\nIn the mouse, CD27 was found on hematopoietic stem cells, multipotent progenitors and common lymphoid precursors (Medina et al., 2001, Nat Immunol 2:718-24; Wiesmann et al., 2000, Immunity 12:193-9).\nCD27 was originally defined as a human T-cell co-stimulatory molecule that increments the proliferative response to TCR stimulation (van Lier et al., 1987, J Immunol 139:1589-96). Presence of CD70 dictates the timing and persistence of CD27-mediated co-stimulation. Upon immune activation, dendritic cells, T-, B- and NK cells transiently express CD70, contingent upon the presence of antigen, Toll-like receptor agonists or inflammatory cytokines.\nCD27 stimulation using anti-CD27 mAb CLB-CD27/1 (9F4) incremented the proliferative response of human T cells to TCR stimulation (Van Lier et al., 1987, J Immunol 139:1589-96). This was confirmed using crosslinked anti-CD27 mAb 1A4, or transfectants expressing CD70. Conversely, antibodies directed to CD27 or CD70 could block this proliferation. Both CD4+ and CD8+ T cells responded to CD27 co-stimulation (Goodwin et al., 1993, Cell 73:447-56; Kobata et al., 1994, J Immunol 153:5422-32; Hintzen et al., 1995, J Immunol 154:2612-23). Studies in mice unambiguously support the role of CD27 as a co-stimulatory receptor for naive CD8+ and CD4+ αβ T cells. For mouse T cells, CD27 primarily promotes their survival upon TCR-mediated activation, but in human T cells, it additionally promotes cell cycle entry and/or activity (reviewed in Borst et al., 2005, Curr Op Immunol 17:275-281; Nolte et al., 2009, Immunol Rev 229:216-231).\nUpon its transient engagement as occurs in acute infections that temporarily upregulate CD70, CD27 supports the generation of a CD8+ effector T cell pool in priming organs, its maintenance at the tissue effector site, its conversion into memory cells and its potential to exercise memory function (Hendriks et al., 2003, J Exp Med 198:1369-1380; Hendriks et al. 2005, J Immunol 175:1666-75, Xiao et al, 2008, J Immunol 181:1071-82). Studies with CD70 blocking antibody in mouse models support the concept that CD27-CD70 interactions can make an important contribution to generation of CD8+ effector T cells, e.g. after protein immunization, virus infection and allotransplantation (Taraban et al., 2004, J Immunol 173:6542-46; Bullock and Yagita, 2005, J Immunol 174:710-17; Yamada et al., 2005, J Immunol 174:1357-1364; Schildknecht et al., 2007, Eur J Immunol 37:716-28).\nTransgenic expression of CD70 in immature dendritic cells sufficed to convert immunological tolerance to virus or tumors into CD8+ T cell responsiveness upon immunization with MHC class I-restricted peptide in PBS. Likewise, agonistic soluble CD70 promoted the CD8+ T cell response upon such peptide immunization (Rowley et al., 2004, J Immunol 172:6039-6046) and in CD70 transgenic mice, CD+ and CD8+ effector cell formation in response to TCR stimulation was greatly facilitated (Arens et al. 2001, Immunity 15:801-12; Tesselaar et al., 2003, Nat Immunol 4:49-54; Keller et al. 2008, Immunity 29:334-346). In mouse lymphoma models, tumor rejection was improved upon CD70 transgenesis or injection of an activating anti-mouse CD27 antibody (Arens et al., 2003, J Exp Med 199:1595-1605; French et al., 2007, Blood 109:4810-15; Sakanishi and Yagita, 2010, Biochem. Biophys. Res. Comm. 393:829-835; WO 2008/051424).\nGenerally, CD27 expression on lymphoid cells is associated with survival potential. Salient examples come from human adoptive T cell therapies, in cancer and AIDS patients, where long-term persisting T cells were selected for CD27 expression. In addition, CD70 expression on tumor-infiltrating lymphocytes was positively correlated with an anti-tumor immune response, potentially reflecting effector T cell survival within the tumor (Ochsenbein et al., 2004, J Exp Med 200:1407-17; Huang et al., 2006, J Immunol 176:7726-35).\nFor conventional CD+ T cells, CD27 similarly promotes primary and secondary responses (Hendriks et al., 2000, Nat Immunol 1:433-40; Xiao et al, 2008, J Immunol 181:1071-82). Moreover, CD27 co-stimulation favours an IL-12 independent pathway for T helper-1 development and enables CD+ T cells to provide help for memory programming of CD8+ T cells (Soares et al., 2007, J Exp Med 204:1095-106; Xiao et al, 2008, J Immunol 181:1071-8). In C57BL/6 mice, CD27 stimulation is consistently associated with Th1-type CD+ T cell differentiation (Arens et al. 2001, Immunity 15:801-12; Soares et al., 2007, J Exp Med 204:1095-106; Xiao et al, 2008, J Immunol 181:1071-82) and in human CD+ T cells in vitro, CD27 promoted Th1 development in presence of IL-12, but had no differentiation-inducing effect in presence of IL-4 (van Oosterwijk et al., 2007, Int Immunol 19:713-18).\nIn addition, CD27 stimulation was demonstrated to promote human regulatory T cell generation and/or function (Jacquot et al., 1997, Cell Immunol 179:48-54). Amongst natural regulatory T cells in human, high CD27 expression hallmarks the cells that have the highest suppressive activity and the CD27 high subpopulation is preferentially amplified during rapamycin treatment (Koenen et al., 2005, J Immunol 174:7573-83; Coenen et al., 2006, Blood 107:1018-23). Recent observations suggest that the CD27+ Treg subpopulation can differentiate into Th17 cells (Koenen et al., 2008, Blood 112:2340-52). Interestingly, CD70+ B lymphoma cells were found to stimulate Treg formation and impede Th17 differentiation by CD27 triggering on intratumoral T cells (Yang et al., 2007, Blood 110:2537-44; Yang et al., 2009, Cancer Res 69:5522-30).\nIn resting B cells, CD27 expression 1s absent, but it is induced during B cell activation in germinal centers and in human, it is subsequently maintained on memory B cells and plasma cells (Agematsu et al., 2000 Immunol Today 21:204-206; Jung et al., 2000, Eur J Immunol 30:2437-2443). CD27 also acts as a co-stimulatory receptor on B cells. In in vitro systems with human B cells, CD27-CD70 interactions consistently stimulate Ig secretion (Agematsu et al., 1997, Eur J Immunol 27:2073-79; Jacquot et al., 1997, J Immunol 159:2652-7).\nHuman NK cells can be subdivided into two functional subsets based on CD27 expression with lack of CD27 expression identifying the mature effector cells (Sugita et al., 1992, J Immunol 149:1199-203; Vossen et al., 2008, J Immunol 180:3739-45). Data suggest a similar co-stimulatory role for CD27 in NK cells as for T-cells (Takeda et al., 2000, J Immunol 164; 1741-1745). The functional effect of CD27 activation on NK cells was established by increased NK mediated killing of CD70-expressing tumor cells. CD27-mediated NK cell activation also promoted the generation of CD8+ anti-tumor immunity (Aulwurm et al., 2006, Int J Cancer 11S:1728-1735; Kelly et al., 2002, Nat Immunol 3:83-90). Recently, NKT cells were shown to promote CD8+ T cell immunity by induction of CD70 on dendritic cells (Taraban et al., 2008, J Immunol 139:1589-96).\nIn addition, CD27 is highly expressed on tumor cells derived from non-Hodgkin's lymphomas and chronic lymphocytic leukemias (Ranheim et al., 1995, Blood 85:3556-3565; van Oers et al., 1993, Blood 82:3430-3436). Soluble CD27 is used as a serum marker for lymphoid malignancy (Van Oers et al., 1993, Blood 82:3430-6).\nIn the research that led to the present invention it was found that the hCD27.15 mAb stimulates the proliferation and/or survival of CD27+ cells. Enhanced proliferation and/or survival of CD27+ cells forms the basis of different therapeutic uses. Monoclonal antibodies that activate CD27 are known. Two activating anti-human CD27 antibodies have been described (Van Lier et al., 1987, J. Immunol. 1987, 139:1589-96; Kobata et al., 1994, J. Immunol. 153:5422-5432). In addition, activating anti-mouse CD27 antibodies have been described (French et al., 2007, Blood 109:4810-15; WO 2008/051424; Sakanishi and Yagita, 2010, Biochem. Biophys. Res. Comm. 393:829-835). hCD27.15 is a unique anti-human antibody, which is, in contrast to 1A4 and 9F4 able to activate human CD27 more effectively than its ligand CD70. These characteristics of hCD27.15 result in a significantly increased effect on proliferation of CD8+ and CD+ T-cells as compared to 1A4, 9F4 and Fc-CD70. Administration of hCD27.15 alone or in combination with other agents to a human being can for example be used in the treatment of cancer.\nCitation or identification of any document in this application is not an admission that such document is available as prior art to the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for hydrodesulfurization of sulfur-containing diesel gas oil, which comprises a specific combination of specific hydrogenation steps.\n2. Description of Related Art\nThe straight run diesel gas oil obtained by distilling the crude oil or the cracked diesel gas oil obtained by cracking heavy oil contains sulfur compounds, and the amount of the compounds is in a range of 1 to 3 wt % as sulfur. When the diesel gas oil containing sulfur compounds is used as a diesel fuel, sulfur compounds will be exhausted into atmosphere as SOx and the environment will be polluted.\nTherefore, these diesel gas oils are used as a fuel usually after being hydrodesulfurized to remove sulfur compounds. It is stated that the permissible value for the amount of sulfur included in a diesel fuel should be 0.05 wt % or less in JIS (Japanese Industrial Standard), and large-scale desulfurization arrangements have been constructed and used to achieve this value. In addition, it is said that it is necessary to decrease the amount of sulfur further with the view of installing a purification catalyst, which reduces NOx in an automotive exhaust gas, into a diesel car in the future, for using a part of the automotive exhaust gas again, by circulating it, as a part of a diesel fuel. This system is called as EGR system (EGR: Exhaust Gas Recirculation).\nA catalyst which consists essentially of cobalt or nickel, and molybdenum, supported on a porous carrier containing alumina as a main ingredient, has conventionally been used for the desulfurization of diesel gas oil so far. However, the conventional catalysts have such problems in that they can hardly remove 4-methyldibenzothiophene and 4,6-dimethyldibenzothiophene, and it is necessary to raise the reaction temperature and the reaction pressure to a very high level in order to lower the sulfur content of the product to the level of 0.05 wt % or less, so that the construction costs of the arrangement and the drive costs increase.\nAs for a process for improving the desulfurization activity to sulfur compounds, a catalyst whose carrier contains phosphorous and boron has been reported in Japanese Unexamined Patent Publication (Kokai) No. 52-13503, as well as a catalyst to whose carrier zeolite is added, has been reported in Japanese Unexamined Publication (Kokai) No. 7-197039. These catalysts have Br.o slashed.nsted acid sites and, thus, exhibit high ability to isomerize a methyl group of dimethyldibenzothiophene and to hydrogenate a phenyl group thereof, and high activity to desulfurize 4-methyldibenzothiophene and 4,6-dimethyldibenzothiophene.\nHowever, catalysts whose carriers contain phosphorous, boron or zeolite, have drawbacks in that their desulfurization activities for alkylbenzothiophenes and dibenzothiophenes without 4- or 6-alkyl substituent, such as dibenzothiophene, and 1-, 2- and 3-methyldibenzothiophene are inferior to those of conventional catalysts consisting essentially of cobalt and molybdenum on alumina carrier (F. van Looij et al. Applied Catalysis A: General 170, 1-12 (1998)). Moreover, said catalysts have further drawbacks in that, as they have Br.o slashed.nsted acid sites, they may easily cause a coloring of the product and when they are used for an olefin-containing feedstock oil or are used at a high temperature of 350.degree. C. or higher, thiols and sulfides are occasionally generated to decrease the desulfurization ratio. In addition, they have another problem in that olefin elements in a feedstock may be polymerized at Br.o slashed.nsted acid sites to generate coke and the deactivation of catalyst may be accelerated. Even if an olefin is not included in a feedstock, if sulfur compounds are desulfurized with said catalyst, an olefin will be generated in situ, and it will cause an extraction of coke. This is understandable from the view that a coking speed, when thiophene flows into said catalyst, reaches ten times the coking speed when olefins or aromatic compounds flow into the catalyst (Catalysis Review, 24, (3), 343 (1982)).\nIt is difficult to desulfurize a diesel gas oil to 0.05 wt % or less as sulfur even if any above-mentioned catalyst is used, and studies have been carried out to deeply desulfurize it from an aspect of a process or reaction apparatus. For example, a process, comprising two different steps under different reaction conditions, which can deeply desulfurize a diesel gas oil without any worsening of hue is proposed in Japanese Unexamined Patent Publication (Kokai) No. 7-102266. A deep hydrodesulfurization process, where a diesel gas oil is separated by distillation into light fraction to be easily desulfurized and heavy fraction to be hardly desulfurized and then these fractions, after being individually hydrodesulfurized, are mixed again into a deep desulfurized diesel gas oil product, is proposed in Japanese Unexamined Patent Publication (Kokai) No. 5-311179. However, said deep hydrodesulfurization process comprising two different steps under different reaction conditions, which can deeply desulfurize a diesel gas oil without any worsening of hue, has an effect to improve a diesel gas oil hue but can hardly improve a further deep desulfurization. Said deep hydrodesulfurization process, where a diesel gas oil is separated by distillation into light fraction to be easily desulfurized and heavy fraction to be hardly desulfurized and then these fractions, after being individually hydrodesulfurized, are mixed into a deep desulfurized diesel gas oil product, has many problems in that a high reaction temperature and a high reaction pressure are needed for heavy fraction to be hardly desulfurized, and the like.\nThus, these prior arts have many problems and they do not achieve an effective production of excellent diesel gas oil with a low sulfur content when used for deep hydrodesulfurization of diesel gas oil as they are."}
{"text": "The prior art discloses various air fresheners including U.S. Pat. No. 4,544,592 to Spector which discloses an aroma generating capsule adhesively attached to the surface of an incandescent bulb. The capsule includes a plastic skin and a vent hole through which an aroma is forced into the atmosphere.\nOf particular interest is U.S. Pat. No. 4,849,606 to Martens which is the same general area of art as the present invention. The Martens patent discloses a tamper-resistant container suitable for use with an electrically heated vapor dispensing apparatus. The container includes a packaged gel which is inserted into a dispensing apparatus having an inner wall with an integrally molded electrical plug and an outer wall having a plurality of openings through which the scented aroma is dispensed.\nU.S. Pat. No. 4,571,485 discloses an electrically heated aromatic generator utilizing a replaceable aroma cube which is inserted into a well in the generator. U.S. Pat. No. 4,731,520 discloses an aroma diffuser assembly for dispensing a scent from a replaceable and expendable cartridge. A heater assembly is also provided to vaporize the cartridge liquid which exits through an aperture into the atmosphere.\nThe prior art also includes U.S. Pat. No. 4,837,421 which shows a fragrance dispensing apparatus which releases a fragrance from a solid polyamide resin. A heating resistor is provided to maintain an elevated temperature. Also of interest are U.S. Pat. No. 3,895,298 to Moran, U.S. Pat. No. 4,425,302 to Pons Pons and U.S. Pat. No. 4,804,821 to Glucksman.\nThe prior art does not disclose an adjustable air freshener of the type set forth herein which is inexpensive and easy to manufacture and use. In the present invention the aroma dispensing can be regulated to conform to the size of a room and waste can be minimized. The disadvantage of an overpowering scent present in the prior art is thus avoided. Further, the useful life of the gel is extended since all of the product is not continuously exposed to the heater at the same time in contrast to the Martens patent.\nCertain factors influence the amount of fragrance released by an air freshener utilizing a heated gel including:\n1. Temperature of the heating element or surface; PA1 2. Surface area of contact between the gel and heated surface; PA1 3. Time of contact; PA1 4 Surface area of permeation (surface area of the film covering the gel); PA1 5. Permeability of the film encapsulating the gel; and PA1 6. The amount of the gel or oil fragrance, and PA1 7. The flow of hot air.\nIn the present invention, factors 1, 4, 5, and 6 are fixed and factor 3 (time) starts from the moment of plug-in and continues until the gel is exhausted. The only factor over which control may be exerted without employing expensive regulating mechanisms is factor 2 (area of contact) and factor 7. The time factor can be regulated by an on-off switch, which is subject to human error, or by an automatic device set at a predetermined interval--both costly and affecting the size of the device. This invention proposes to regulate the scent from an air freshener or dispenser by controlling the area of contact between the heater and gel with a rotatable cover which cooperates with a unique base. The dial is rotated to place the scented gel in contact with the heated surface to produce the desired scent. A dividing wall across the center of the cover restricts the hot air affecting only the portion of gel which is moved up to the heated compartment. The portion of gel which remains in the cold compartment is protected from the contact of the hot air by the dividing wall."}
{"text": "1. Field of the Invention\nThe present invention relates allowing a memory controller to optimize strobe to clock relationship for a DRAM.\n2. Background Information\nA Dynamic Random Access Memory, DRAM, is a typical memory to store information for computers and computing systems, such as, personal digital assistants and cellular phones. DRAMs contain a memory cell array having a plurality of individual memory cells; each memory cell is coupled to one of a plurality of sense amplifiers, bit lines, and word lines. The memory cell array is arranged as a matrix of rows and columns, and the matrix is further subdivided into a number of banks.\nOne type of DRAM is a synchronous dynamic random access memory (SDRAM) that may allow for synchronous operation with a processor. Specific types of SDRAM are a single data rate (SDR) SDRAM and a double data rate (DDR) SDRAM. For example, a present DDR memory type is DDR3. Typically, a DDR DRAM may send data (DQ), when enabled by a DQS clock signal, to a memory controller or memory controller hub (MCH). The memory controller or MCH may receive the data from the DDR DRAM by utilizing precision delay cells to provide a delayed DQS clock signal.\nTypically, motherboard specification and guidelines synchronize the DQS to a system clock (SCLK) relationship at a specific DRAM. However, with the advent of DDR3 DRAM, a typical tree structure for clocks is not utilized and the DQS to SCLK relationship differs by DRAM position on a Dual Inline Memory Module (DIMM)."}
{"text": "Storage devices such as disks are well known. It was also well known that storage controllers control access to the storage devices in response to read and write requests. The storage controllers also mirror data to different storage devices and spread data amongst different storage devices for redundancy and backup purposes. Storage controllers were also known to store data in accordance with one of several known RAID security levels. Generally, the higher the RAID level the greater the redundancy of the data storage.\nIt was also known that pooled storage devices can be used to increase storage capacity and provide recovery and backup services.\nStorage servers, such as an IBM Enterprise Storage Server, were also known. An IBM ESS storage server includes two clusters of processors and associated hardware. Typically, there are four storage controllers in each cluster. Each of the storage controllers controls multiple storage devices grouped in RAID arrays. In one environment, clients with Fiber Channel Host Bus Adapters (“HBAs”) are coupled via a Fiber Channel to a switch. The switch is also coupled to the Storage Server with Fiber Channel HBAs. There may be multiple storage servers per client. Each client is assigned or allocated storage “volumes” which are mapped to physical locations on storage devices that are grouped in RAID arrays. Consequently, clients make data access requests (reads and writes) to the storage server, for data within their allocated volumes, and the storage server accesses the mapped locations in cache storage to satisfy the requests or from disk if the data does not reside in cache storage.\nA known IBM Enterprise Storage Server comprises two clusters of storage controllers with four storage controllers per cluster. Each cluster has its own cache (semiconductor) memory shared by all storage controllers in the cluster. Each cluster also has battery backed up nonvolatile storage (“NVS”) which is shared by all of the storage controllers in the cluster, although any one storage controller cannot use more than 25 percent of the total NVS allocated to the cluster. In addition, each storage controller has its own cache (disk) storage. The cache memory is used for rapid access to data inpaged from external storage to service read data access requests from memory and to provide buffering of modified data. All update requests are written to the associated cluster and backed up by the NVS on the mirrored cluster.\nOccasionally, one or more of the storage controllers or arrays become “over-utilized”, i.e. there are more data access requests than can be handled expeditiously. The problem may be due to I/O delays within the storage server or other hardware or bandwidth constraints. It was known to move some of the data from a disk of an over-utilized or over-accessed array to another, lesser-utilized array. It was also known that when a cache memory is too small, there may be excessive outpaging of the data in the cache memory to storage, to make room for inpaging of data for which access is requested. It was known to increase the size of the cache memory, when found to be too small, to the capacity of the hardware.\nA known IBM TotalStorage Expert program tool monitored cache memory, NVS, Physical Array Activity and Logical Volume activity in a known IBM Enterprise Storage Server. The TotalStorage Expert program tool is further described in published documents, “IBM TotalStorage Expert Reporting: How to Produce Built-In and Customized Reports” by Daniel Demer et al., dated October 2003, and “IBM TotalStorage Expert Hands-On Usage Guide” by Mary Lovelace et al., dated January 2003.\nA known Tivoli Multiple Device Manager program tool monitored cache memory, NVS, Physical Array Activity and Logical Volume Activity in a known IBM Enterprise Storage Server, and reported performance data and exceptions. The Tivoli Multiple Device Manager program tool is further described in a published document, “Managing Disk Subsystems using IBM TotalStorage Productivity” by Mary Lovelace et al., dated September 2005.\nAn object of the present invention is to accurately determine the effectiveness of a storage system to access a storage device.\nAnother object of the present invention is to take corrective action when the effectiveness of the storage system is impaired."}
{"text": "1. Field of the Invention\nThis invention relates to apparatus for electric wiring, and more particularly to an apparatus for pulling an electric wire through an electric conduit.\n2. Description of the Prior Art\nWhen a building is constructed, electric conduits are provided in each house for electric wiring. When an electric wire is to be installed in an electric conduit, a spring wire is used and inserted from one end of the electric conduit to the other end thereof, and then the electric wire to be installed is fastened to the lead end of the spring wire, and then the spring wire is pulled back to carry the electric wire through the electric conduit. This electric wiring method consumes much labor and working time."}
{"text": "Refractory metals such as tantalum and titanium have been used as diffusion barriers, adhesion layers or X-ray lithography masks in electronic applications. Unfortunately, these materials oxidize easily when annealed causing reliability problems due to the increased resistivity and stress. Other refractory metals such as Nb, V, Zr, Hf, Cr, Mo and W also oxidize easily in annealing ambients containing minimal amounts of oxygen and at relatively low temperatures. It is an expensive and difficult process to remove all oxygen from an annealing environment and some processes actually require an oxidizing ambient. Small amounts of oxygen incorporated interstitially in the refractory metals can cause large increases in compressive stress. For instance, in a Ta film that has 10% oxygen incorporation, an increase of 4 to 5 GPa in compressive stress and an increase of 63 .mu..OMEGA.-cm in resistivity have been observed. When used in electronic applications, such increases in resistivity and stress can cause major reliability problems such as poor contact resistance, peeling and cracking.\nOther researchers have utilized refractory metals in electronic applications as diffusion barriers, i.e., Ta diffusion barrier between Si and Cu was used by Holloway, et al., J. Appl. Phys., 71 (11), 5433 (1992). The use of Ta diffusion barrier between a high dielectric material and Si in the fabrication of a high dielectric oxide material has been shown in Grill, et al., J. Mater. Res., 7 (12), 3260 (1992). In the latter work, the fabrication of high dielectric constant oxide materials requires high temperature processing (&gt;650.degree. C.) in oxygen ambients which generally cause the underlying Si to oxidize, and thus creating an additional in-series low dielectric capacitor.\nOther diffusion barrier materials such as TiN, WN and TaN have been used in VLSI applications, i.e., in contact holes for separating conductors from underlying silicon. However, these materials are not suitable oxygen diffusion barriers because they cannot withstand oxidation anneal cycles that the devices may be subjected to. In addition, other researchers have investigated elemental metals such as Pt, Au and Ru for the prevention of diffusion of oxygen to the underlying layer of silicon and its subsequent oxidation. It was found that none of the pure metals prevented the diffusion and the resulting SiO.sub.2 formation. A break in the electrical conduction path to the silicon substrate occurred as a result.\nIn certain semiconductor processes, high dielectric constant Perovskite compounds (such as PZT, PLZT or Ba.sub.x Sr.sub.1-x TiO.sub.3) are deposited onto a substrate. These materials require high temperature (&gt;650.degree. C.) oxygen anneal in order to crystallize. It also requires a Pt seed layer for crystallization into a Perovskite phase which has the highest dielectric constant, i.e., 380 vs. about 40. For simple process integration, it is desirable to use Si as the lower electrode. A problem thus incurred during an oxygen anneal of the Perovskite compound/Pt/Si structure is the formation of an SiO.sub.2 layer at the Pt/Si interface which reduces the effective dielectric constant. A diffusion barrier layer is therefore needed for preventing oxygen diffusion down to the Si layer, and also for preventing the Si diffusion up to the Pt layer.\nIt is therefore an object of the present invention to provide a diffusion barrier for use in a semiconductor device that does not have the shortcomings of the prior art diffusion barriers.\nIt is another object of the present invention to provide a diffusion barrier that can be fabricated in a simple manufacturing process.\nIt is a further object of the present invention to provide a diffusion barrier that is capable of sustaining the high processing temperatures encountered in semiconductor processing steps.\nIt is another further object of the present invention to provide an oxygen diffusion barrier for devices incorporating high dielectric constant oxide layers.\nIt is yet another object of the present invention to provide an oxygen diffusion barrier of a multi-layer aluminum/refractory metal material such that the refractory metal layer does not oxidize (remains electrically conducting) in a subsequent annealing step.\nIt is still another object of the present invention to provide an oxygen diffusion barrier that is a multi-layer thin film of aluminum/refractory metal such that when the film is exposed to atmosphere a thin layer of Al.sub.2 O.sub.3 forms on the surface of aluminum."}
{"text": "Collapsible strollers, especially child strollers, of the known type typically include a front-mounted safety bar which is collapsible to enhance portability. During use of known strollers, a stroller occupant could become injured by being pinched between the movable sections of the collapsible bar.\nIt is desirable to provide a stroller with a collapsible tray having a reduced pinching hazard."}
{"text": "The reference in this specification to any prior publication (or information derived from it), or to any matter which is known, is not, and should not be taken as an acknowledgment or admission or any form of suggestion that that prior publication (or information derived from it) or known matter forms part of the common general knowledge in the field of endeavour to which this specification relates.\nIdentification of new compounds for use in pharmaceutical preparations is an important part of the search for more reliable and effective therapies. However just as important is the development and modification of known compounds, reducing the risks associated with a new drug candidate and significantly reducing the development and cost to bring the drug to clinical development.\nMany drugs fail in clinical trials either because their physical properties (particularly solubility) make them difficult to formulate, or because of a poor therapeutic index that leads to toxic effects during the high drug concentrations that occur just after dosing. Other short comings include poor absorption, poor bioavailability, instability, systemic side effects due to an inability to target the drugs, and the inability to control their biodistribution, metabolism and renal or hepatic clearance once administered. Similarly some current products on the market can be improved with regard to such issues.\nA number of approaches have been tried to improve a pharmaceutical compound's profile including the formulation of the pharmaceutical agent in a liposome, micellar or polymeric micelle formulation, as well as covalent attachment of the pharmaceutical agent to a hydrophilic polymer backbone.\nThe characteristics of an ideal profile modifying agent include being a well defined structure, allowing precise control of the absorption, distribution, metabolism and excretion (ADME) characteristics (also referred to as pharmacokinetics) of the compound in question and advantageously being able to carry multiple compounds per agent or construct. The toxicity of a compound in question can be ameliorated through its controlled release from the said agent or construct, the body only being exposed to therapeutic plasma concentrations of the compound.\nIn recent years, dendritic macromolecules, or dendrimers, have been found to have increasing applications in biotechnology and pharmaceutical applications. Dendritic macromolecules are a special class of polymers with densely branched structures that are characterized by higher concentrations of functional moieties per unit of molecular volume than ordinary polymers. There are three subclasses of dendritic macromolecules: random hyperbranched polymers; dendrigraft polymers and dendrimers (which includes dendrons), classified on the basis of the relative degree of structural control present in each of the dendritic architectures (Fréchet and Tomalia “Dendrimers and other Dendritic Polymers”, Wiley and Sons, New York, 2002). The unique properties of dendrimers in particular, such as their high degree of branching, multivalency, globular architecture and well-defined molecular weight, make them promising new scaffolds for pharmaceutical applications. In the past decade, research has increased on the design and synthesis of biocompatible dendrimers and their application to many areas of bioscience including macromolecular drugs, drug delivery, biomedical imaging and medical devices.\nThe potential utility of dendritic polymers both as drug delivery vectors and pharmaceutical actives has received increasing interest in recent years. However, whilst the literature is replete with reports of, for example, synthetic schemes for dendrimer assembly, descriptions of dendrimer-drug interactions and drug loading efficiencies and increasingly, in vitro evaluations of dendrimer interactions with cell lines, there is very little information describing the fundamental pharmacokinetic and metabolic fate of dendrimers.\nFurther, it is still a challenge to prepare dendrimers that circulate in the blood long enough to accumulate at target sites, but that can also be eliminated from the body at a reasonable rate to avoid long-term build up. In addition, the tissue localisation of dendrimers is still difficult to predict in advance and more studies are required to determine the effect of peripheral dendritic groups on these properties. An additional area that needs to be investigated is the release of drugs from dendrimers. Steric hindrance associated with the dense globular dendritic architecture makes the engineering of the enzymatically cleavable linkages difficult.\nIt has surprisingly been found that, by introducing one or more functional moieties, as described herein, the efficacy of the macromolecule may be significantly improved, but without significant adverse impact on, or interference with, other functional moieties which may be present."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Unless otherwise indicated herein, the materials described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nMachine-Type Communication (MTC) is a promising and emerging technology. Potential MTC based applications include smart metering, healthcare monitoring, remote security surveillance, intelligent transportation systems, etc. These services and applications stimulate the design and development of a new type of MTC device that may need to be integrated into current and next generation mobile broadband networks such as LTE and LTE-Advanced.\nThe existing mobile broadband networks were designed to optimize performance mainly for human type of communications and thus are not optimized for MTC specific requirements. For instance, some MTC devices may be installed in basements of residential buildings. These devices may experience significantly greater penetration losses on the radio interface than normal LTE devices. In order to provide sufficient coverage of such MTC devices, special coverage enhancement considerations may be needed."}
{"text": "Various hydraulic, hydraulic/pneumatic, and fully electrical systems with all-electrical actuators are known for active vertical adjustment of a motor vehicle. An active “spring strut” that is suitable for this vertical adjustment has a suspension spring and a preloaded spring which are each separated by an adjustable spring cap. The adjustable spring cap is exposed to positioning forces during vertical adjustment. That is, the adjustable spring cap as needed traverses an axial path dictated by the control, and thus ultimately produces a positioning force on the wheel. The lifting motion of the adjustable spring cap takes place almost exclusively by a ball groove thread drive in which the adjustable spring cap does not rotate during its lifting motion, but remains torsionally strong, while the relatively slender spindle of the ball groove thread drive rotates and is driven by an electric motor which is located in the spring strut itself.\nThis type of wheel suspension is described by DE 10 2005 001 744 B3 in which there is a preloaded spring located concentrically on the spring strut above the suspension spring; this preloaded spring opposes the spring force of the support spring and thus reduces the adjusting forces which must be applied by an electromechanical actuator or a positioning drive in order to effect adjustment of the level and/or pitching and rolling equalization on the motor vehicle by way of the axially adjustable spring cap. The three-dimensional and functionally reliable arrangement of the positioning drive integrated on the spring strut within the preloaded spring is a problem, however, especially for wheel suspensions with steered wheels in which the wheel requires a corresponding free space in the wheel well and for the wheel suspension elements.\nThe spindle shown in DE 10 2005 001 744 B3 does have a large bearing base. Valuable installation space is lost, however, by lining up the elongated bearing bush with its lower bearing arrangement. This has an adverse effect on the shock absorber length and on the shock absorber stroke. Moreover, the total weight of the spring strut is increased by the long piston rod which is required in this design and the long bearing bush. The shock absorber plunges into the lower part of the spindle and is cardanically deflected relative to the spindle when driving. This can lead to contact between the oscillating shock absorber and the rotating spindle. This contact between the shock absorber and spindle is further promoted by the piston rod which is elastically supported in its upper end. In order to prevent this contact between the shock absorber and spindle, the radial air gap between the shock absorber and spindle would have to be dimensioned to be correspondingly large. This would, however, lead to a relatively large spindle diameter, as a result of which, however, the rotational mass moment of inertia of the spindle is adversely increased.\nThe object of the invention is to propose a spring strut arrangement of the generic type which ensures a structurally especially compact and functionally reliable construction."}
{"text": "An example of a conventional engine support structure is described in Unexamined Japanese Patent Application Publication No. 2002-160663, Vehicle Body Front Portion Structure. For the structure described in Unexamined Japanese Patent Application Publication No. 2002-160663, as shown in FIG. 4 (a), sub-frame 51 is disposed below vehicle body front frame member 50, which extends in the front to rear direction of the vehicle, so that a predetermined distance is formed above and below said vehicle body front frame member 50 and so that it extends in the front to rear direction of the vehicle, and the front end of aforementioned vehicle body front frame member 50 and the front end of sub-frame 51 are connected by means of connecting member 52, which extends in the upper to lower direction. In addition, the engine is loaded onto aforementioned sub-frame 51.\nTherefore, when a predetermined load or more is input from the front of the vehicle toward the rear due to a collision, a means is devised to increase the energy absorbing efficiency of this input load by collapsing and deforming the aforementioned vehicle body front frame member 50.\nConventionally, as cited in Unexamined Japanese Patent Application Publication No. 2003-246276, for example, the vehicle body and sub-frame are connected via a supporting elastic body, such as an insulator for absorbing the vibration, and a spacer for filling the space that is formed between the vehicle body and the sub-frame.\nFor the technology disclosed in aforementioned Unexamined Japanese Patent Application Publication No. 2003-246276, a spacer is disposed between the vehicle body and the sub-frame into which an insulator has been inserted. This spacer and insulator are securely tightened to the vehicle body by means of a bolt that penetrates in the upper to lower direction through the bottom of the sub-frame. In addition, below the insulator and sub-frame is disposed a stopper for preventing the sub-frame from falling out due to breakage caused by deterioration of the insulator over time."}
{"text": "Certain computer operating systems and applications are designed to load modules that are internally requested by the operating system, internally requested by an application, or externally requested by a supported process. Such modules may include, for example, a procedure, a function, a script, a plug-in module, a driver, or a portion of an executable application.\nIn some computer operating systems and/or applications, modules that can be loaded are classified according to a module type. For example, in an operating system, the module type may be used to determine how to properly support the module. In some cases, a sub-system of the operating system, such as a dynamically loadable kernel module (“DLKM”) sub-system, for example, may be used to support modules of different types.\nFIG. 1 illustrates a prior operating system configuration 100 for handling dynamically loadable modules. A table 110 of defined module types and associated type-specific support module identifiers are included in a statically configured operating system kernel 105. Support module identifiers may refer to support modules, such as support module 120, that contain, for example, pre-loading logic, post-loading logic, pre-registration logic, post-registration logic and registration logic.\nWhen the operating system 105 receives a request to load a module 115, the operating system uses the module type definition table 110 to determine whether the module type is supported. Upon locating the correct module type entry in the table, the operating system utilizes the support modules (e.g., support module 120) identified in the table 110 to properly support the requested module 115.\nThe operating system 105 may provide module type-specific support for only the module types defined by the table within the statically configured kernel. This type of operating system configuration is unable to dynamically load modules of types that have not been previously defined in the statically configured portion of the operating system kernel.\nAccordingly, whenever a new module type is to be defined, the statically configured kernel must be reconfigured to include a definition of the new module type. The statically configured kernel must be reloaded to accommodate the newly added module type definition. The reconfiguration of the statically configured kernel results in an increased kernel size and a reloading or rebooting of the operating system."}
{"text": "1. Field of the Invention\nThe invention relates to a driving method of a plasma display panel (hereinafter, simply referred to as a xe2x80x9cPDPxe2x80x9d) of a matrix display type.\n2. Description of Related Art\nAs one of the display panels of the matrix display type, a PDP of AC (alternating current discharge) type is known.\nThe AC type PDP has a plurality of column electrodes (address electrodes) and a plurality of row electrode pairs which are arranged so as to perpendicularly intersect the column electrodes and in which one scanning line is formed by one pair. Each of the row electrode pairs and the column electrodes are covered by a dielectric layer against a discharge space and a discharge cell corresponding to one pixel is formed at an intersection point of the row electrode pair and the column electrode.\nAs a method of allowing the PDP to embody a halftone display, there is a so called a subfield method wherein one field period is divided into N subfields each emits the light only for a time corresponding to a weight of each bit digit of pixel data of N bits and the data is displayed. The method is disclosed, for example, in Japanese Patent Kokai No. 4-195087.\nFIG. 1 is a diagram showing a format of the light-emission driving in one field period according to the subfield method.\nIn the example shown in FIG. 1, in which it is presumed that the pixel data being supplied consists of six bits, the period of time of one field is divided into six subfields SF1, SF2, . . . , and SF6, to perform the light-emission driving. By executing the light emission through the six subfields once, an expression of 64 gradations can be performed for an image of one field.\nEach subfield is constructed by an all-resetting step Rc, a pixel data writing step Wc, and a sustain light emitting step Ic. In the all-resetting step Rc, by discharge-exciting (reset discharge) all of the discharge cells of the PDP in a lump, wall charges are uniformly formed in all of the discharge cells. In the next pixel data writing step Wc, a selective erasure discharge according to the pixel data is excited every discharge cell. The wall charges in the discharge cell in which the erasure discharge has been performed are extinguished and this discharge cell becomes a xe2x80x9cnon-light emission cellxe2x80x9d. The discharge cell in which the erasure discharge is not performed becomes a xe2x80x9clight emission cellxe2x80x9d because the wall charges remain there. In the sustain light emitting step Ic, a discharge light emitting state is continued only for the light emission cell only for a time corresponding to the weight of each subfield. In each of the subfields SF1 to SF6, consequently, a sustain light emission is sequentially performed with weights of light emitting period ratios of 1:2:4:8:16:32.\nFor example, if an image such that a flat object moves is displayed by the driving method, however, there is a problem such that a fringe-like pseudo outline like a video image as if the gradations were lost is observed at points near the positions where the luminance gradation level transverses the boundary of n-th power of 2, such as xe2x80x9c32xe2x80x9d or xe2x80x9c16xe2x80x9d.\nThe above problem occurs because in the case where the luminance gradation level is xe2x80x9c32xe2x80x9d, the light emission is performed only in the subfield SF6 during the 1-field period as shown in FIG. 1 and, when the luminance gradation level is equal to xe2x80x9c31xe2x80x9d, the light emission is not performed in the subfield SF6, and the light emission is performed in the subfields SF1 to SF5. That is, for a period of time during which the discharge cell which should perform the light emission at the luminance gradation level xe2x80x9c32xe2x80x9d is lit on, the discharge cell which should perform the light emission at the luminance gradation level xe2x80x9c31xe2x80x9d is certainly in a light-off state, so that the fringe-like outline that is not concerned with the image is confirmed on the boundary of the discharge cells.\nTo suppress the pseudo outline and improve a display quality, therefore, a driving method whereby a subfield of a relatively long light emitting period is further divided into a plurality of subfields and they are distributed and arranged in the 1-field period has been proposed. According to the driving method, as a light emission pattern in the 1-field period is uniformed by increasing the number of subfields, a suppressing effect of the pseudo outline rises.\nThere is, however, a problem such that as the number of subfields increases, the number of bits of drive data that is formed in correspondence to each subfield has to be also increased in accordance with it and a scale of the apparatus increases with the increase in number of bits.\nThe invention is made to solve the problems and it is an object of the invention to provide a driving method of a plasma display panel, in which a scale of an apparatus can be reduced by suppressing the number of bits of drive data while maintaining an image display of a high quality in which a pseudo outline has been suppressed.\nAccording to the invention, there is provided a driving method of a plasma display panel in which a display period of one field is divided into a plurality of subfields and a light emitting state in each of the subfields is set in accordance with each bit of pixel data, to perform a halftone display, wherein the light emitting state in a subfield of a relatively long light emitting period is also set in accordance with the bit of the pixel data which sets the light emitting state in a subfield of a relatively short light emitting period among the subfields."}
{"text": "This invention relates to variable capacity wobble plate compressors mainly adapted for use in air conditioning systems for automotive vehicles, and more particularly to an improved wobble plate compressor in which the crankcase pressure is controllable for varying the displacement or capacity of the compressor.\nA variable capacity wobble plate compressor in general is adapted to vary its displacement or capacity through a change in the angularity of the wobble plate. It is known e.g. from U.S. Pat. No. 3,861,829 to vary the refrigerant pressure in the crankcase for changing the angularity or angle of inclination of the wobble plate relative to the drive shaft. A conventional wobble plate compressor of this type comprises a fluidtight housing, a drive shaft rotatably disposed in the housing, a cylinder block arranged in the housing and formed therein with a plurality of cylinders circumferentially arranged around the drive shaft and extending substantially parallel to the axis of the drive shaft, pistons received in the respective cylinders for reciprocating motions therein, and a wobble plate supported at its diametrically central portion by trunnion pins extending at right angles to the drive shaft and axially movable therealong and also supported at its peripheral edge by a pivot pin rotatable about the drive shaft together therewith. The wobble plate is adapted to be pivotally displaced in unison with axial movement of the trunnion pins along the drive shaft to have its angularity varied relative to the drive shaft. As the wobble plate rotates in a position inclined relative to the drive shaft, the pistons are reciprocatingly moved in their respective cylinders for pumping actions. In the compressor, the resultant reaction force exerted by all the pistons, some on their compression strokes and some on their suction strokes, acts upon the wobble plate at a point inside a half portion of the circumference described by the axes of the cylinders, which is located at the same side of the drive shaft as the pistons on their compression strokes, so that the wobble plate is acted upon by the above resultant reaction force to become inclined relative to the drive shaft about the trunnion pins as a movable fulcrum during the pumping actions of the pistons. The resultant reaction force of the pistons counteracts the pressure in the crankcase which acts upon the pistons as back pressure. Therefore, when there occurs a drop in the pressure in the crankcase, the wobble plate is displaced in the angularity-increasing direction to increase the capacity, whilst when there occurs an increase in the crankcase pressure, the wobble plate is displaced in the angularity-decreasing direction to decrease the capacity.\nIn this arrangement, it is known that the difference between values of the crankcase pressure corresponding, respectively, to the maximum angularity of the wobble plate and the minimum angularity thereof falls within a range from 5 to 10 percent of the difference between the suction pressure and the discharge pressure of the compressor. For example, if the compressor is operating in a condition wherein the discharge pressure is 14 kg/cm.sup.2, and the suction pressure is 2.1 kg/cm.sup.2, the crankcase pressure has to be controlled within a very small range from approximately 2.7 kg/cm.sup.2 to approximately 3.3 kg/cm.sup.2, with a small pressure difference of approximately 0.6 kg/cm.sup.2. This means that the angularity of the wobble plate varies in high response to a change in the crankcase pressure. This requires precise control of the crankcase pressure, making it difficult to control the capacity in a stable manner."}
{"text": "Integrated optical devices (e.g., photonic integrated circuits) are well suited to applications in various technologies such as telecommunications, instrumentation, signal processing and sensors. In operation, photonic integrated circuits use optical waveguides to implement devices, such as optical switches, optical couplers and wavelength multiplexers/demultiplexers, for example. Such waveguides, when integrated with a photonic integrated circuit, are typically implemented as solid dielectric light conductors, which are fabricated on a substrate in a very similar fashion as semiconductor integrated circuits are manufactured. Waveguides transmit light around optical circuits and also connect to external optical waveguides, such as optical fibers, typically by direct physical abutment of the fiber with the waveguide. However, in such a configuration, a mode mismatch results between the integrated waveguide and the optical fiber.\nAs such, improvements are needed."}
{"text": "High density optical recording systems which may be used for recording and playing back information are well known in the art. Current practice in the optical recording field can be broadly distinguished into three categories, the classification being largely dependent on the physical and chemical characteristics of the information carrying media. Common to all three categories is the ability of each system to have the information stored on the media to at least be \"read\" or \"played back\" by irradiating selected areas of the media with a highly focused source of light, usually a laser. Changes in the transmission or reflection of the laser beam caused by the presence or absence of recorded information on the media is detected by suitable optical sensors. The three categories are differentiated by the media's response to radiation. \"Read Only\" as, for example, in Compact Disk technology, wherein all the information on the disk is pre-recorded. \"Write-Once, Read (Mainly)\" (WORM), in which information can be written on the media by the laser, read back by a laser at reduced power levels, but cannot be erased for correction or re-use. The third consists of media which can be written upon and read back by laser, but can also be erased (by the laser), in the presence of a suitable magnetic field, and which then can be reused.\nThis invention, relates generally to the \"erasable\" optical recording systems, known in the art as magneto-optical recording systems, and in which the media is described as magneto-optic media.\nA magneto-optic recording medium is a magnetic material which causes the polarization angle of laser light to be changed when it is reflected from or transmitted through a recorded spot. Preparatory to the initial recording operation, the magneto-optic media is vertically magnetized in the same direction across the entire recording surface. During recording, the direction of magnetization of the media is reversed at the point of incidence of the laser beam due to the presence of a vertical magnetic field of a given strength and polarity and heating induced by a laser. Illumination from a playback laser beam of lower power reflected from this area will show a polarization angle change, which is detectable by suitable optical sensors. The erasing process restores the original magnetic direction which existed prior to recording by the use of a vertical magnetic field of polarity opposite to that used for recording when accompanied by heating from the laser beam. Traditionally, an electromagnet in which the current can be adjusted and reversed, (thereby changing the direction and/or intensity of the induced magnetic field) has been the preferred source for generating the magnetic bias for both the writing and the erasing. The electromagnet (in some cases, a reversible permanent magnet) is also sometimes referred to as the bias magnet.\nDeguchi et al, in U.S. Pat. No. 4,497,006, schematically describe in their disclosure an electromagnet surrounding the objective lens of their recording/playback apparatus which provides the bias magnetization required for magneto-optic operations. Ota et al, in U.S. Pat. No. 4,477,852, discloses a pair of bias electromagnets adjacent the objective lens in conjunction with a pair of lasers for performing the erase and re-write functions in a magneto-optic recording and reproducing system. U.S. Pat. No. 4,466,004 also describes an electromagnet surrounding the objective lens, and records information by varying the field strength of the electromagnet in the areas of the media heated by a continuous laser beam. In the previous two examples, writing was accomplished by varying the intensity of the laser beam.\nU.S. Pat. No. 4,610,009, assigned to the Xerox Corporation, describes a system wherein plural lasers operate through a single optical path and cooperate with bi-polar magnetic field producing means located on the opposite side of the magneto-optic medium. This system also includes a bias electromagnet essentially concentric with the focusing electromagnet of the objective lens. The deficiencies of these examples in providing an essentially vertical magnetic field coincident with the laser beam is addressed in U.S. Pat. No. 4,340,914 issued to Hanaoka and assigned to Olympus Optical. This patent discloses an auxiliary magnetic pole piece for focusing a magnetic flux from the magnetic field generating device (bias magnet) into at least the heated portion and which is located on the side of the media opposite the laser beam.\nThe teachings of U.S. Pat. No. 4,701,894, assigned to RCA, further address the problem of providing a vertical magnetic field at the point of incidence (with the laser beam) which rely on an electromagnet located around the objective lens, or, in some cases, on the side of the media opposite the laser beam. As discussed in this patent, given the very close proximity between the objective lens of the record/playback optics, there is very little space on the recording side of the media for an electromagnet capable of providing a vertical magnetic field at the point of incidence. This patent discloses a bias electromagnet configured in the shape of a letter \"E\", in which the top and bottom legs of the \"E\" extend radially from one edge of the media towards the center (when the media is in a disk configuration, the center leg of the \"E\" likewise extends from one edge towards the center, but on the opposite side of the media from the top and bottom legs). Current carrying coils are wound about each of the two segments of the vertical portion of the \"E\" to magnetize the electromagnet. Attached to the housing of the objective lens is a ring or cap of magnetic material which has a small hole in it to allow passage of the laser beam. Magnetic coupling of this cap to the top and bottom legs of the \"E\" is accomplished by a pair of pole pieces, in effect, filling in the gap between the legs and the cap. The pole pieces, the cap, and the objective lens can move radially in-between the top and bottom of the \"E\", with the return path of the magnetic flux being provided by the center leg which is located underneath the lens and the media.\nWhen the magneto-optic media is provided on a thin tape or on a thin disk, as described in this patent, the magnetic elements can conveniently be placed on the side of the media opposite the side upon which the laser impinges. However, when this opposite side location is not physically accessible, for example, when the magneto-optic media is supported on the surface of a drum, or is supported on a very thick substrate, these magnetic elements must be located on the same side of the media as the optics are.\nIt is an object of this invention to provide a device and method of producing the necessary vertical bias magnetic field, essentially in concurrence with the point of incidence of the laser beam upon the media, in which access to the opposite side of the media is denied. It is a further object of this invention to provide the necessary vertical bias magnetic field for systems in which extensive cross-track movement of the optic system is required. It is an additional object of this invention to provide the necessary vertical magnetic field with a reduced number of non-vertical components of the bias magnetic field, and with greatly reduced power requirements for providing the bias magnetic field."}
{"text": "1. Field\nThe present disclosure relates generally to light-emitting diode (LED) bulbs, and more specifically to using a laminate structure for mounting LEDS in a liquid-filled LED bulb.\n2. Description of Related Art\nTraditionally, lighting has been generated using fluorescent and incandescent light bulbs. While both types of light bulbs have been reliably used, each suffers from certain drawbacks. For instance, incandescent bulbs tend to be inefficient, using only 2-3% of their power to produce light, while the remaining 97-98% of their power is lost as heat. Fluorescent bulbs, while more efficient than incandescent bulbs, do not produce the same warm light as that generated by incandescent bulbs. Additionally, there are health and environmental concerns regarding the mercury contained in fluorescent bulbs.\nThus, an alternative light source is desired. One such alternative is a bulb utilizing an LED. An LED comprises a semiconductor junction that emits light due to an electrical current flowing through the junction. Compared to a traditional incandescent bulb, an LED bulb is capable of producing more light using the same amount of power. Additionally, the operational life of an LED bulb is orders of magnitude longer than that of an incandescent bulb, for example, 10,000-100,000 hours as opposed to 1,000-2,000 hours.\nWhile there are many advantages to using an LED bulb rather than an incandescent or fluorescent bulb, LEDs have a number of drawbacks that have prevented them from being as widely adopted as incandescent and fluorescent replacements. One drawback is that an LED, being a semiconductor, generally cannot be allowed to get hotter than approximately 120° C. As an example, A-type LED bulbs have been limited to very low power (i.e., less than approximately 8 W), producing insufficient illumination for incandescent or fluorescent replacements.\nOne approach to alleviating the heat problem of LED bulbs is to use a thermally conductive liquid to cool the LEDS. To facilitate thermal dissipation, it may be advantageous to increase the thermal paths from the LED to the environment."}
{"text": "1. Field of the Invention\nThe present invention relates generally to strip-line-type circuits, more particularly to loop transmission lines as shunt capacitors used in such circuits.\n2. Description of the Related Art\nCapacitors are one of the basic building blocks for electronic and microwave circuits. In microwave engineering, strip-line-type circuits, including microstrip, strip-line, and multi-layer circuits, can use large metal conductor patches to approximate shunt capacitors. Such patch capacitors can be found, for example, in bias networks of amplifiers, microstrip low pass filters, and matching networks.\nAs with parallel plate capacitors, the capacitance realized from conductive patches on a microstrip circuit is directly proportional to the area of the patches and the dielectric constant of the substrate. Examples of microstrip patch capacitors are shown in FIG. 1. These patch capacitors can occupy a significant amount of surface area, depending on the amount of capacitance required and the type of the substrate used. Use of conductive patch shunt capacitors thus places significant limitations on the layout flexibility and minimum sizes of circuits.\nIt is thus desirable to construct shunt capacitors that offer more layout options or potential for more compact circuit design or both. The present invention is directed to achieve one or more of these goals.\nIn accordance with the principles of the invention, a strip-line-type circuit includes a shunt capacitor that includes a closed conductive loop. The circuit may further include a transmission line connected to the closed conductive loop. The transmission line may be connected to the closed conductive loop at two nodes, in which case the closed conductive loop is divided into two segments, connected in parallel at the two nodes. The impedance of one of the two segments may be substantially larger than the impedance of the other segment, as in the case, for example, where one segment is substantially longer than the other.\nThe closed conductive loop may be a layer of conductive thin-film pattern formed on a layer of dielectric material, including a loop made of a superconductor such as YBa2Cu3O7-d (YBCO) formed on a magnesium oxide, sapphire or lanthanum aluminate substrate.\nThe circuit may be a multi-layer circuit in which the closed conductive loop extends to multiple layers of conductive patterns.\nThe closed loop may take on a variety of shapes, including circular, rectangular and swirl shapes.\nMore particularly, the circuit may be a filter that includes an inductor with each of its ends connected to a closed conductive loop that acts as a shunt capacitor. The filter may include multiple inductors connected in series, with the junctions between the inductors connected to shunt capacitors realized by closed conductive loops.\nThe filter may be constructed from a variety of materials, including the above-listed examples. For example, the filter may a band-stop filter having five or more poles constructed from YBCO film on a magnesium oxide substrate no larger than about 50 mm in any dimension."}
{"text": "Ambrisentan, a potent selective endothelin receptor antagonist, is utilized for the treatment of pulmonary arterial hypertension and idiopathic pulmonary fibrosis. It blocks the endothelin receptor on vascular smooth muscle cells and cardiac myocytes, preventing vasoconstriction and smooth muscle proliferation.\nAmbrisentan is marketed under the trade name Letairis™ or Volibris®. Ambrisentan is chemically known as (S)-2-((4,6-dimethylpyrimidin-2-yl)oxy)-3-methoxy-3,3-diphenyl propanoic acid and has the structural formula as shown in formula (I).\n\nU.S. Pat. No. 5,703,017 broadly discloses several molecules generally, which are 3-(Het) aryl carboxylic acid derivatives. The compounds disclosed also include Ambrisentan. In the specification, synthetic schemes for preparation of the molecules and for their intermediates are disclosed. The general methods of preparation involve coupling the epoxide of formula III (for example, with R═COOR) with alcohols or thiols of formula IV according to the following scheme:\n\nU.S. Pat. No. 7,109,205 discloses and claims Ambrisentan. The patent also discloses the process for the preparation of Ambrisentan as per a similar scheme as in the above patent. This specification also discloses the preparation of Ambrisentan according to the above process.\nThis patent also describes the preparation in laboratory scale of (S)-2-hydroxy-3-methoxy-3,3-diphenylpropanoic acid (compound V) by resolution of the corresponding racemate using L-proline methyl ester and (S)-1-(4-nitrophenyl)ethylamine. By this resolution, yield of 35% based on the racemate and an optical purity of 99.8% were achieved.\nU.S. Pat. No. 6,559,338 describe a process of resolution of 2-hydroxypropanoic acid by using an optically active base (S)-1-(4-chlorophenyl)ethylamine and then one of the diastereomeric salts formed is separated off. However the optically active base used here is very expensive and would not be preferable for large-scale production. Additionally, the objective of this resolution has not been clearly defined.\nHowever, the U.S. Pat. No. 6,559,338 also state that, it has been found that when this described reaction step was scaled up (several kg to 100 kg), additional working steps became necessary in order to ensure a high optical purity. The diastereomeric salt of (S)-2-hydroxypropanoic acid and (S)-1-(4-nitrophenyl)ethylamine crystallizes with difficulty and therefore cannot be filtered off readily either, so that some of the mother liquor remains in the crystals together with the enantiomer to be separated off. Only when the crystals were additionally stirred in the tank together with fresh solvent, and when the crystals that had been filtered off are copiously rewashed, the required optical purity was obtained.\nTherefore, the present invention provides a process for resolution of racemic 2-hydroxy-3-methoxy-3,3-diphenylpropanoic acid (VIII) that does not have the above-mentioned disadvantages but can readily be carried out on an industrial scale.\nIn addition, the inventors of the present invention surprisingly found the improved resolution of the racemate compound (VIII) by reacting the racemic acid with a suitable optically active amine and subsequently separating off the diastereomeric salt with high chiral purity. Additionally the same diastereomeric salt can further be utilized for the preparation of Ambrisentan.\nThus, the present invention discloses an improved process for preparing Ambrisentan. The present invention discloses an improved process for preparing Ambrisentan wherein the chiral amine base used as resolving agent, provides a novel intermediate that is used in preparation of Ambrisentan of high enantiomeric purity.\nWO2010017918 disclosed the amorphous form of Ambrisentan and process for the preparation of thereof. It is well known that amorphous forms have several limitations in terms of stability, processability etc. thereby making the crystalline form of any compound more preferable. The present inventors have surprisingly found a new crystalline form of Ambrisentan which is stable, easily processable & can be used to obtain the commercial form of Ambrisentan in chirally and chemically pure form."}
{"text": "1. Technical Field\nThe present invention relates to a recording apparatus such as an ink jet type printer, and a medium loading device provided in the recording apparatus.\n2. Related Art\nGenerally, as a kind of recording apparatus, an ink jet type printer is widely known. This kind of printer performs printing by supplying ink into a recording head and ejecting the supplied ink from a nozzle of the recording head toward a recording medium. Some printers of this kind use large-size roll paper as the recording medium (see, for example, JP-A-2009-23171).\nIn the printer according to JP-A-2009-23171, for preparation of printing, a roll assembly formed by integrating a roll of a long recording medium that is rolled into a roll, a pair of flanges attached to both end sections of the roll, and a spindle member inserted into the center of the roll is loaded into a loading portion.\nIn the printer according to JP-A-2009-23171, for loading the roll assembly into the loading portion, the roll assembly is temporarily placed on a temporary placement portion, and is subsequently moved by being rotated on a surface of a guide portion so that the roll assembly may be transferred to the loading portion. Thus, there is a problem that, when loading the roll assembly into the loading portion, the roll assembly may become unstable.\nSimilar problems may occur not only in the above-described ink jet type printer but also in all kinds of recording apparatuses that use a roll assembly."}
{"text": "1. Field of the Invention\nDisclosed herein is a flashing beacon. More specifically, disclosed herein is a flashing beacon that may be positioned on most any roadway sign or signpost, and that may include a signal unit, a control unit, a solar collector, an activation device (e.g., a timer, microwave emitter, radio transmitter, step-pad, a pushbutton, infrared transmitter, wireless transmitter or like device) and various other accessories.\n2. Reference to Related Art\nAccording to the U.S. Manual on Uniform Traffic Control Devices, flasher mechanisms associated with traffic control signs (e.g., a yield or crosswalk sign) must be positioned on the sign (or signpost) so that flashing signal is about 12 feet above the pavement. The flashing signal must also be programmed or otherwise set to flash continuously at a rate of not less than 50 nor more than 60 times per minute. See MUTCD, Section 4D.11. However, while the guidelines set forth in the uniform regulations are intended to provide a visible warning to drivers, recent testing has suggested that only a small percentage of the public responds to flashing signals that operate according to the uniform regulations. Specifically, recent testing has suggest that only 25% of the public complies with or otherwise responds to flashing signals associated with roadway signage. Therefore, it would be advantageous to have an improved flashing beacon system that may be used with existing or future roadway signage to garner a greater response from the vehicle driving public."}
{"text": "1. Field of the Invention\nThe present invention pertains to apparatus and methods for obtaining core samples from subterranean formations. More specifically, it pertains to apparatus and methods for obtaining core samples from the bottom of a well bore, primarily in producing and exploring for oil and gas or other minerals. In particular, the present invention is directed to apparatus and methods for obtaining pressurized core samples and raising them to the surface of the well without losing liquids and gases contained therein and without creating an overpressured and potentially dangerous core barrel.\n2. Brief Description of the Prior Art\nIt has long been a practice in drilling of oil and gas wells to monitor or examine the subterranean formations in which the well is being drilled by obtaining a core sample at the bottom of the well bore. Customarily, a core is obtained by drilling an annulus at the bottom of the well bore, leaving a cylindrical core section extending upwardly within a core barrel of some type. The core section or sample is then broken off and retrieved to the surface for examination and analysis. In the past, such coring was done by cable tool devices or by removing the drill string with the core sample captured therein. The core sample is of course initially at the bottom of the well where it may be under very high formation pressures. In core devices of the prior art, the natural pressure within the core is reduced as the core is brought to the surface until it is essentially at atmospheric pressure. As this pressure is reduced, it is probable that much of the liquids and gas, which are naturally captured in the core, will be released so that once the core reaches the surface, it is difficult to determine what may have been contained in the strata at the bottom of the hole.\nBecause of the problems associated with losing formation pressure, core barrel devices have been developed which capture the core at the bottom of the well, then seal it within a pressurized barrel for removal to the surface. One such cable device is shown in U.S. Pat. No. 2,347,726. Other pressure capturing core barrel devices for attachment to the lower end of a drill string are shown in U.S. Pat. Nos. 3,064,742; 3,146,837; and 3,454,117. However, since the pressure at the bottom of a well hole may be extremely high, e.g. 10,000 p.s.i., the core in such an assembly may be encapsulated in a cylinder which is potentially a bomb when it reaches the surface. It may be hazardous to disassemble and requires extremely careful handling at the surface with specialized equipment.\nIn recent years, one or two core sampling devices have been developed which totally capture a pressurized sample, but allows for some reduction of pressure within the core barrel as it is raised to the surface so that when the core sample reaches the surface, it is at a much more practical pressure for working with, yet retains the liquids and gases present in the sampled strata. One such device is disclosed in U.S. Pat. No. 2,287,909 (Sewell). In Sewell, a rubber sleeve, silfon bellows, rubber baloon or other type of expansion member is provided in fluid communication with the chamber in which the core sample is received. However, the type of expansion members disclosed in Sewell are of relatively small volume and are not suitable for the high pressure formations encountered in modern drilling practice. It is not uncommon to encounter formations in which gas or vapors therein will expand several hundred times at atmospheric pressure. The coring device of Sewell could not accommodate such expansion and would probably fail well before reaching the surface.\nWhile othbrs may have attempted to develop coring apparatus suitable for capturing high pressure core samples for sealed retrieval to the surface, while allowing the pressure within the device to be reduced as it is brought to the surface of the well. Since such sampling apparatus and methods are highly desired for proper sampling and analysis, the search will undoubtedly continue to overcome the disadvantages of coring apparatus and methods of the prior art."}
{"text": "Internet access has advanced from the early dialup system to take advantage of improving wired and wireless technologies. High capacity data networks are currently offered over cable, fiber connections, and wireless networks. For example, cellular systems operate 3 G and 4 G networks that utilize new and efficient protocols, such as Worldwide Interoperability for Microwave Access (WiMAX) and Wireless Metropolitan Area Networks (WMAN) among others to provide increasing bandwidth and coverage.\nThese improvements in throughput are, however, being matched by increasing demand for services such as audio and video streaming, Internet protocol (IP) television, and games that use the IP to provide significant volumes of data at the highest rates possible. The current standards for multimedia streaming, such as Microsoft Adaptive/Smooth streaming and Apple's HTTP Live Streaming, require additional power and bandwidth for each stream.\nThe challenge for IP service providers is to deliver as many services and to support as many customers as possible at a competitive price. The price to the subscriber of an IP service is directly related to the investment in infrastructure needed to provide a satisfactory and competitive service.\nWired network access devices may be configured to access bandwidth that is provided by the wired service provider and wireless bandwidth that is provided by a wireless service provider or that is more generally accessible and relatively inexpensive or free (a wired device with this capability may be referred to herein as a “wired multi-network device”). Wireless devices are increasingly able to access bandwidth that is provided by the wireless service provider and other wireless bandwidth that is relatively inexpensive or free (a wireless device with this capability may be referred to herein as a “wireless multi-network device”).\nThe most prevalent protocol currently used to provide the non-service provider wireless connectivity is “WiFi.” A multi-network device (wired or wireless) may communicate over one or more service provider networks (wired or wireless) and over wireless access points that may be located in homes, offices or at so-called “hotspots.” When connecting to an access point, the multi-network device is in an “infrastructure mode.”\nMulti-network devices may also have the ability to communicate directly with each other or with a group of devices. This mode of communication is referred to as “ad-hoc” or “peer-to-peer” mode. For example, a group of WiFi devices may form a “mesh” network or a “peer-to-peer” network. In this mode, there is no master base station or access point.\nFor wireless multi-network devices, the WiFi capability provides an alternative network connection when the device is near a WiFi access point (AP) to off-load the mobile network traffic from the relatively more expensive wireless service provider network. However, the off-loading of network traffic from the wireless service provider does little to improve the overall throughput of the service provider network in the face of demands for large data files."}
{"text": "The applicant's U.S. Pat. No. 4,399,778 describes a method of operating a two-stroke or two cycle high torque and high speed gasoline or diesel fueled multi-cylinder internal combustion engine utilizing an air blower or turbo charger combination connected in series, and having a new and improved design of the internal moving parts, more specifically the scavenging valves which are opened and closed by the wide area of cam lobes of a cam shaft and related mechanisms all assembled in the cylinder head. Also of importance is the piston assembly having a novel feature of design. A recirculating type constant pressure lubrication system is used to lubricate all the vital and critical moving parts, making this invention an efficient device."}
{"text": "Conventionally, barcodes have been used as measures for providing paper or plastic media with specific information and recognizing the information automatically. For commodity management and logistics, information of barcodes provided on commodities is read by using a barcode reader such as a laser scanner and a CCD camera. Coming up in these years is a new type of barcode reader that scans a whole surface of an ID card for capturing an image and recognizes a barcode based on the image, while a personal number being recorded as the barcode on the ID card. For this type of barcode reader, it is important to accurately detect a barcode symbol from a 2D image space. Some detection methods are known for the detecting process (Refer to Patent Documents 1 and 2).\nIn a process of Patent Document 1, by using a degree of similarity of derivatives of scanning lines parallel in a local region of an image space as an evaluation value, it is judged whether the region includes a barcode image or not. The derivatives are derived by differential operation. Meanwhile, in a process of Patent Document 2, a differential waveform of a candidate area possibly including a barcode is calculated for quantifying some feature quantities included in the waveform, and it is judged according to the quantified values whether the region includes a barcode image or not. Also in the latter case, the differential waveform is calculated by differential operation.\n[Patent Document 1]\nJapanese Unexamined Patent Application Publication No. 4-225487\n[Patent Document 2]\nJapanese Unexamined Patent Application Publication No. 2004-252897\nUnfortunately, the processes disclosed by Patent Documents 1 and 2 require complicated calculations including differential operations for identifying a position of a barcode symbol, as described above, so that a lot of computations are needed. That may result in an obstacle to process time shortening. Meanwhile, a high-performance arithmetic element can be used for shortening the process time. However, such an arithmetic element is costly, and may lead to an increased cost of the entire reader as a result."}
{"text": "In many applications, such as radar warning receivers (RWRs), there is a need to detect short-pulse signals with high probability.\nFrequency-swept detectors and spectrum monitors are not well-suited to being short-pulse signal detectors since they suffer low dwell times in each frequency band.\nShort-pulse signals can be detected using a digital instantaneous frequency measurement (IFM) technique based on an analogue frequency discriminator and reference is made to P. L. Herselman and J. E. Cilliers: “A Digital Instantaneous Frequency Measurement Technique using High-Speed Analogue-to-Digital Converters and Field Programmable Gate Arrays” (2007). A conventional frequency discriminator typically comprises a splitter, a delay line, a mixer and a low-pass filter. Using the splitter, delay line and mixer, an input signal is mixed with a delayed copy. The product is fed through a low-pass filter. The frequency of a single input tone can be found using a look-up table. However, one drawback of such a system is that it can “blinded” by an extra input tone. Another drawback is that such a system tends to suffer from having a limited dynamic range.\nUS 2011/0053538 A1 describes a high-speed frequency sensor. The sensor comprises a plurality of filters allowing passage of signals over a different band of frequencies and a plurality of detectors, each associated with a respective filter. If the incident RF signal results in a signal in the band of frequencies passed by a particular filter, the detector associated with that filter generates an output signal. This indicates in which of the frequency bands the incident RF signal has been detected. The precision of the frequency estimate is limited to one bin per filter, that is, the sensor determines only a range of frequencies within which the single frequency exists, the range being equal to the band of frequencies passed by the filter."}
{"text": "Laundry washing machines are used in many single-family and multi-family residential applications to clean clothes and other fabric items. Due to the wide variety of items that may need to be cleaned by a laundry washing machine, many laundry washing machines provide a wide variety of user-configurable settings to control various aspects of a wash cycle such as water temperatures and/or amounts, agitation, soaking, rinsing, spinning, etc. The settings cycle can have an appreciable effect on washing performance, as well as on energy and/or water consumption, so it is generally desirable for the settings used by a laundry washing machine to appropriately match the needs of each load washed by the machine.\nSome laundry washing machines also support user selection of load types, typically based on the types of fabrics and/or items in the load. Some laundry washing machines, for example, have load type settings such as colors, whites, delicates, cottons, permanent press, towels, bedding, heavily soiled items, etc. These manually-selectable load types generally represent specific combinations of settings that are optimized for particular load types so that a user is not required to select individual values for each of the controllable settings of a laundry washing machine.\nWhile manual load type selection in many cases simplifies a user's interaction with a laundry washing machine, such manual selection still can lead to suboptimal performance due to, for example, user inattentiveness or lack of understanding. Therefore, a significant need continues to exist in the art for a manner of optimizing the performance of a laundry washing machine for different types of loads, as well as reducing the burden on users when interacting with a laundry washing machine."}
{"text": "Certain entities, such as large corporations, government agencies, or universities, may maintain a backbone network (also referred to as a “computer network”) to interconnect entity assets. For example, a corporate backbone network may be used to connect different data centers maintained by a corporation. Backbone networks may additionally include one or more points of presence for connecting the backbone network to the Internet. Backbones networks are themselves made up of connectivity devices, such as hubs, switches, and routers, as well as cables to connect the connectivity devices. These components of the backbone network provide the network capacity by which the network demand between assets (e.g., from data center to data center) may be satisfied.\nTo increase backbone network capacity, an entity's network planners may engage in both short-term operational planning and long-term strategic planning. In the short-term, network planners may, for example, acquire additional capacity from an existing cable in the backbone network that has unused capacity, such as by paying for increased cable capacity. In the long-term, network planners may utilize new cables in the backbone network, thereby adding new routes between data centers in the network or augmenting existing routes between data centers.\nTraditionally, the long-term planning done by network planners has been driven by the availability of new supply. That is, as new cables have been laid, network planners have had to decide whether to negotiate for use of the new cables, either through leasing access to the cables or joining a consortium of entities that manage the cable. These decisions have typically been made by network planners based on the current network demands of their backbone networks. In other words, the decisions behind long-term strategic planning have been made by network planners based primarily on the state of the backbone network at that particular time. Such limited input into long-term decision making is inherently inefficient, and risks the ability of the backbone network to meet future network demands. The currently available tools for network planning are inefficient. The planning of network can include planning a transport layer of the network, e.g., physical link layer such as laying of fiber optics, their capacities, and an Internet Protocol (IP) layer, e.g., placement of routers in the network, how the routers have to be connected. The current tools typically facilitate planning one layer or the other and do not support multi-layer planning. It would therefore be beneficial to facilitate strategic planning of a backbone network."}
{"text": "1. Field of the Invention\nThe present invention relates to an antenna device capable of transmitting and receiving radio signals in a plurality of frequency bands, and a communication terminal apparatus including such an antenna device.\n2. Description of the Related Art\nIn a communication terminal apparatus including a mobile phone, for example, such a loop antenna as disclosed in Japanese Patent Laying-Open No. 2002-43826 may be utilized. This loop antenna is configured by a looped-shaped conductor having one end as a power feed end and the other end as a ground end, and having an entire length of one wavelength. This loop antenna suppresses gain reduction even when being used in proximity to a human body, and exhibits excellent radiation characteristics.\nIn recent years, there is a need for a communication terminal apparatus to accommodate a plurality of frequency bands. For example, a communication terminal apparatus accommodating a penta-band of GSM (registered trademark; Global System for Mobile communication) 850, GSM900, GSM1800, GSM1900, and UMTS (Universal Mobile Telecommunications System) is required to accommodate a relatively wider band of 824 to 960 MHz (Low Band) and 1710 to 2170 MHz (High Band).\nAccording to the loop antenna for accommodating such a relatively wider band, as shown in FIG. 1A, three resonances (resonance 1, resonance 2 and resonance 3) are used to cover a plurality of frequency bands. In other words, resonance 1 forms a passband in a Low Band while resonance 2 and resonance 3 form a band in a High Band.\nAs shown in FIG. 1B, resonance 1 is caused by fundamental waves in the odd mode, and shows a resonance mode having monopole-type current distribution in which the intermediate point of loop antenna 101 is defined as an electric field maximum point. Resonance 2 occurs in the even mode, and shows a resonance mode having dipole-type current distribution in which there are two electric field maximum points on loop antenna 101. Resonance 3 is caused by harmonics in the odd mode, and shows a resonance mode having current distribution as shown in the figure in which there are three electric field maximum points on loop antenna 101. In this case, the “odd mode” represents a mode in the state where the current direction from the power feed end to the radiation element and the current direction from the ground end to the radiation element are aligned with each other. The “even mode” represents a mode in the state where the current direction from the power feed end to the radiation element and the current direction from the ground end to the radiation element are opposite to each other.\nThe resonance frequency of each resonance can be determined by the size of loop antenna 101. On the other hand, when this resonance frequency is controlled in a matching circuit, it is conceivable to implement a configuration in which an inductance element L1 and an inductance element L2 are loaded at the power feed end and the ground end, respectively, of the antenna, as shown in FIG. 1C.\nHowever, when inductance elements are loaded in this way to adjust the frequency, the amount of change in each resonance frequency is increased as the frequency is higher. In other words, by the method of simply loading an inductance element, it is difficult to independently control the resonance frequency for each resonance mode."}
{"text": "This invention relates to a device (hereinafter referred to as a “support adaptor assembly”) which can be fitted to a pressing device (hereinafter referred to as a “bearing press”) used for removing or for fitting press-fitted components such as bearings, gears, pulleys, pins, bushes and seals.\nConventional bearing presses comprise a support frame having a base section, two spaced apart elongate side members extending upwards from the base section, and a fixed mounting member to which is connected a mechanically or hydraulically operated pressing device. Attached to the elongate side members of the support frame by pins or bolts, are two parallel, horizontal, spaced apart support members used to support a component to which pressure is applied by operation of the pressing device.\nThere are many variations of this configuration, but all known prior art bearing presses have two parallel, horizontal, spaced apart support members, above which is mounted a pressing device.\nIn use, an object requiring assembly or dismantling of press-fitted components, is placed on the support members of the bearing press, or on plates supported by the support members of the bearing press, so that the pressing device can be operated to apply vertical, downward pressure on the object to remove or replace the press fitted components. In order to prevent damage to components and injury to the operator of the bearing press, the component to which pressure is applied must be on a horizontal plane relative to the vertical direction of the downward pressing force.\nThere are, however, some objects with press-fitted components, which cannot be correctly aligned with a horizontal plane or supported on the support members of a bearing press for the following reasons:                a) the object may be too long to fit between the pressing device and the horizontal support members of the bearing press. Even when the support members have a means of height adjustment, this is often a difficult and time consuming process;        b) the object may be too wide to fit between the horizontal support members of the bearing press;        c) the object may be of an irregular shape which does not allow it to be supported in such a way that places the component to be fitted or removed with a horizontal plane relative to the vertical direction of the downward pressing force;        d) the object may comprise two components of similar width or diameter which are press fitted together and must be separated by use of a pressing force for repairs.        "}
{"text": "The present invention relates generally to methods and systems for performing flow cytometry, and more particularly to such methods and systems that allow performing in-vivo flow cytometry in a live subject without extracting blood.\nMany current techniques for detecting and quantifying various cell types circulating within a subject require extracting a blood sample from the subject followed by analysis of the extracted blood. For example, one such conventional method includes withdrawing a blood sample, fluorescently labeling specific cell populations ex vivo with antibody-targeted fluorescent markers, and passing those cells through a flow stream to be interrogated by a light source (typically a laser). Fluorescence and light scattering signals emitted, or reemitted, by the cells in response to the interrogating light can be employed to determine the types and the number of cells. Automated hematology analyzers can separate and count blood cell populations based on their size/volume characteristics in addition to their light scattering properties. Another blood analysis method, known as hemocytometry, is based on counting cells against a grid while being viewed with a microscope.\nThese techniques, though useful, require invasive blood withdrawal and hence are not particularly suited for continuous, real-time monitoring of the circulation. Further, blood withdrawal can be difficult in certain patient populations, e.g., children and psychiatric patients. In addition, such blood withdrawal can expose patients to various risks, such as infection and/or anemia. For example, in patients with low blood volume (e.g., newborns and premature infants), blood loss due to phlebotomy for common tests can equal as much as 45% of the total blood volume. Hence, diagnostic blood withdrawal is among the leading causes of anemia in neonates, which can lead to the need for frequent transfusion or erythropoietin treatment. Further, for monitoring patients for potential infection or in diseases such as leukemia and other types of cancer, frequent monitoring of the white blood cell count (WBC) can be beneficial but is currently not practical due to the frequency of blood withdrawal required. Additionally, blood withdrawal can be dangerous in immune-compromised patients, such as AIDS patients, transplant patients, or patients undergoing chemotherapy, as the skin penetration for extracting blood carries an inherent risk of infection.\nSome methods for performing in-vivo flow cytometry have also been proposed previously. These methods typically utilize fluorescent antibodies or fluorescent proteins, such as green and red fluorescent proteins (e.g., GFP and DsRed), for labeling cells. However, fluorescent antibodies and fluorescent proteins are not approved for human use and require either genetic manipulation or intravenous injection. As such, the applicability of such methods can be limited.\nAccordingly, there is a need for improved methods and systems for performing flow cytometry in vivo, and particularly for such methods and systems that would allow non-invasive, real-time, continuous monitoring of cell populations in circulating blood."}
{"text": "1. Field of the Invention\nThe present invention relates to an inverted pendulum type vehicle capable of traveling on a floor surface.\n2. Description of the Related Art\nThere has conventionally been known an inverted pendulum type vehicle in which a rider mounting section tiltable relative to the vertical direction is attached to a base body, to which a travel operation unit that travels on a floor surface and an actuator that drives the travel operation unit are installed (refer to, for example, Japanese Patent Application Laid-Open No. 2011-068165).\nA rider of the inverted pendulum type vehicle moves the vehicle by moving his/her weight or by operating an operation unit, such as a joystick. It requires, however, a certain level of skill to move the inverted pendulum type vehicle exactly as the rider intends. Hence, it has been difficult for a beginner, who is not experienced with the operation, to operate the inverted pendulum type vehicle as he or she intends."}
{"text": "This specification relates to scoring images that are related to particular entities.\nThe Internet provides access to a wide variety of resources, for example, webpages, images, audio files, and videos. A variety of search engines are available for identifying particular resources accessible over the Internet. For example, digital images that satisfy a user's informational need can be identified by an image search process in which keywords or other data are processed to identify collections of digital images. Each image is identified to the user by an image search result. The image search result may include a thumbnail of an image to which it links.\nSearch results generated by a search engine are generally ranked and presented to a user in an ordered list of search results. The rankings may be based on how the search engine determines the relevance of resources specified by the search results to a particular query issued by the user.\nA very common entity that people search for is a person. Thus, users use search engines to identify images associated with a particular person. The images associated with the particular person typically are images that depict the face of the particular person."}
{"text": "1. Field of the Invention\nThis invention relates generally to measurement and data acquisition systems and, more particularly, to programmable gain instrumentation amplifier (PGIA) design.\n2. Description of the Related Art\nScientists and engineers often use measurement systems to perform a variety of functions, including measurement of a physical phenomena or unit under test (UUT), test and analysis of physical phenomena, process monitoring and control, control of mechanical or electrical machinery, data logging, laboratory research, and analytical chemistry, to name a few examples.\nA typical measurement system comprises a computer system with a measurement device or measurement hardware. The measurement device may be a computer-based instrument, a data acquisition device or board, a programmable logic device (PLD), an actuator, or other type of device for acquiring or generating data. The measurement device may be a card or board plugged into one of the I/O slots of the computer system, or a card or board plugged into a chassis, or an external device. For example, in a common measurement system configuration, the measurement hardware is coupled to the computer system through a PCI bus, PXI (PCI extensions for Instrumentation) bus, a GPIB (General-Purpose Interface Bus), a VXI (VME extensions for Instrumentation) bus, a serial port, parallel port, or Ethernet port of the computer system. Optionally, the measurement system includes signal conditioning devices which receive the field signals and condition the signals to be acquired.\nA measurement system may typically include transducers, sensors, or other detecting means for providing “field” electrical signals representing a process, physical phenomena, equipment being monitored or measured, etc. The field signals are provided to the measurement hardware. In addition, a measurement system may also typically include actuators for generating output signals for stimulating a unit under test.\nMeasurement systems, which may also be generally referred to as data acquisition systems, may include the process of converting a physical phenomenon (such as temperature or pressure) into an electrical signal and measuring the signal in order to extract information. PC-based measurement and data acquisition (DAQ) systems and plug-in boards are used in a wide range of applications in the laboratory, in the field, and on the manufacturing plant floor.\nIn a measurement or data acquisition process, analog signals may be received by a digitizer, which may reside in a DAQ device or instrumentation device. The analog signals may be received from a sensor, converted to digital data (possibly after being conditioned) by an Analog-to-Digital Converter (ADC), and transmitted to a computer system for storage and/or analysis. When a measurement system generates an analog output signal, the computer system may generate digital signals that are provided to one or more digital to analog converters (DACs) in the DAQ device. The DACs convert the digital signal to an analog output signal that is used, e.g., to stimulate a UUT.\nThe input stage of many measurement devices (e.g., a DAQ device) typically consists of one or more instrumentation amplifiers. Generally, analog signals that are received at a DAQ device are first routed from a particular channel via a multiplexer. The signals may then enter an instrumentation amplifier, typically a programmable gain instrumentation amplifier (PGIA). A distinctive feature of an instrumentation amplifier is that it provides very high input impedance and gain. A PGIA typically applies a specified amount of gain to an input signal, which amplifies or attenuates the signal and ensures proper A/D conversion. The PGIA may also convert differential input signals applied to the DAQ board to a single-ended output so that the ADC can correctly digitize the data.\nAs the resolution and speed of ADCs increase, the accuracy, noise, and speed requirements on PGIAs may similarly increase. Some common PGIA designs may not be able to meet the additional requirements. One conventional design is a three op-amp PGIA design. This design may suffer from a decrease in bandwidth as the gain is increased. Its accuracy and linearity may also suffer at high gains due to the reduction of negative feedback necessary to effect higher gain. Also, the rejection of common-mode signals may be limited by resistor matching.\nOther PGIA designs typically suffer from output offset errors and drifts due to matching errors between multiple DC current sources and mirrors. These current sources and mirrors may also contribute a significant amount of output-referred noise. Also, in most PGIA designs, as more components are added to the input stage, the settling time of the overall amplifier may suffer significantly. Each component may introduce additional delays due to parasitics, some of which exhibit dielectric absorption effects, and therefore may change the overall time constants and step response of the amplifier. Moreover, device tolerances, mismatches, and any asymmetrical topology in the PGIA may result in a relatively low common mode rejection ration (CMRR).\nFurthermore, PGIAs may be susceptible to overload-induced errors, such as saturation recovery time or excessive current draw. Insufficiently protected PGIAs may pass along overload conditions to the next stage of circuitry, typically the ADC, which may not be designed to handle such overload conditions."}
{"text": "1. Technical Field of the Invention\nThe present invention relates to the decoding of radiofrequency transmission channels conveying coded digital information.\nThe invention thus applies advantageously to digital terrestrial television, that is to say using signals transmitted between television antennas, as defined, for example, in the European DVB-T (Digital Video Broadcasting-terrestrial) or in the North-American ATSC (Advisory Committee on Advanced Television Service Digital TV) specification, or to digital cable television, as defined, for example, in the European DVB-C (Digital Video Broadcasting-cable) or in the North-American ITU-J83 Annex A/B/C/D (International Telecommunication Union) specification, all these telebroadcasts being based on the MPEG transmission standards, and use for example to convey information, quadrature digital modulation, or else modulation of the COFDM type according to terminology well known to the person skilled in the art.\nThe invention also relates in particular to tuners.\n2. Description of Related Art\nThe television signals received at the input of a tuner of a receiver are composed of the entire set of channels transmitted that lie in the 45 MHz-860 MHz frequency band. The objective of the tuner is to select the desired channel and to output a baseband signal on the in-phase path (I path) and on the quadrature path (Q path). This signal is thereafter converted into a digital signal and demodulated. The channel decoding processing thereafter also comprises a block which distinguishes, typically by means of majority logic, zeros from ones, then performs all the error correction, that is to say typically a Viterbi decoding, the interleaving, Reed-Solomon decoding and deshuffling. The channel decoding device outputs packets that are decoded in a conventional manner in a source decoding device in accordance with the MPEG standards so as to reoutput the initial audio and video signals transmitted via antennas or via the cable.\nCurrent studies are aimed at researching ever more integrated solutions for the embodiment of digital terrestrial or cable television receivers. However, this research hits technological difficulties related to the embodiment of the means making it possible, on the one hand, to correctly discern a channel from the entire set of channels present, and, on the other hand, to correctly discriminate the signal from noise. Specifically, whereas in digital satellite television, all the channels have nearly the same power, this is not the case in particular in digital terrestrial television. Thus, it is for example possible to have adjacent channels exhibiting a very high power with respect to the desired channel, for example a difference in power of the order of 40 dB. It is therefore necessary to be able to ensure rejection of 40 dB on the adjacent channels. Moreover, the specification imposes a signal/noise ratio of the order of 30 dB. Hence, in the worst case, it is therefore essential for it to be possible to bring the level of an adjacent channel to 70 dB below its input level. This results in extremely significant constraints on the filters, this being extremely constraining in respect of integrated solutions.\nThere exists a need to solve this problem. More specifically, there is a need for a totally integrated solution for the tuner, given the characteristics of digital television signals, in particular terrestrial ones, exhibiting more significant rejection constraints than digital cable television signals."}
{"text": "1. Field of the Invention\nThe present invention relates to an adjusting apparatus for automatically adjusting positional deviation of stereoscopic images by means of geometrical image transformation processes.\n2. Discussion of the Background Art\nIn recent years, a stereoscopic vehicle surrounding monitoring apparatus using a pair of left and right cameras (stereoscopic camera having solid image element like CCD) mounted on the vehicle has been interested by automobile engineers. To detect a distance to an object, first respective picture element blocks having coincidence of brightness are found in left and right images (stereo matching), then distance data are calculated according to the principle of triangulation from a parallax, namely a relative deviation amount, between both picture element blocks. Consequently, in order to calculate distance data with high reliability, it is desirable that there exists no positional deviation other than the parallax in a pair of left and right images (stereoscopic images). In actual world, however, the stereoscopic camera has some amount of positional error caused when the camera is installed on the vehicle.\nWith respect to this, Japanese Patent Application Laid-open No. Toku-Kai-Hei 10-307352 discloses a technology in which the positional deviation of the stereoscopic camera is corrected by applying a geometric transformation to the stereoscopic image. That is, when an initial adjustment of the positional deviation is made or when a readustment of the positional deviation generated by aged deterioration is made, a dedicated correction detecting device is connected with an image correction apparatus performing the affine transformation to calculate the difference of angle of view, a rotational deviation or a parallel deviation of the stereoscopic image obtained by imaging a specified pattern for adjustment and to establish (reestablish) parameters of the affine transfomation according to the result of the calculation. The positional deviation is equivalently corrected by applying the affine transformation to images based on thus established affine parameters.\nHowever, according to the prior art, a special adjustment pattern is imaged by the stereoscopic camera and the deviation is corrected based on the position of the pattern in images. Accordingly, when the correction is performed, it is necessary to interrupt the ordinary surroundings monitoring control and as a result this prior art is not suitable for a real time processing in which the monitoring control is carried out concurrently.\nIt is an object of the present invention to provide an adjusting apparatus capable of correcting the positional deviation of stereoscopic images in real time in parallel with the surroundings monitoring control.\nIn order to attain the object, a positional deviation adjusting apparatus of a stereoscopic camera taking a first image or a reference image and a second image or a comparison image, comprises an image correcting means for transforming an image geometrically according to transformation parameters, a reference area establishing means for establishing a plurality of reference areas in a first image transformed by the image correcting means, a reference point establishing means for establishing a reference point in the reference area for each of the reference points, a reference line establishing means for establishing a reference line passing through the reference points, a searching area establishing means for establishing a searching area corresponding to the reference area in a second image, a correlation point identifying means for identifying a correlation point having a correlation with the reference point for each of the searching areas, an approximation line calculating means for calculating an approximation line based on the correlation points and a parameter correcting means for correcting the transformation parameters such that the approximation line coincides with the reference line."}
{"text": "1. Field of the Invention\nThe invention relates to a knitting machine with at least one guide bar, which is moveable by a shogging drive in a shogging direction and by a swing-through drive perpendicular to the shogging direction.\n2. Discussion of Background Information\nIn the production of a knitted article in a knitting machine, guides must be guided around needles in order to produce a loop formation. The guides are thereby attached to the guide bar. The sequence of movements of the guide bar is therefore controlled such that in a loop formation the guide bar is moved once back and once forth in a shogging direction. Thus, the shogging direction is the direction that corresponds to the longitudinal extension of the guide bar, i.e., the direction in which all of the guides of the guide bar are arranged one behind the other. The guide bar is additionally moved perpendicular to the shogging direction between the forth (forward) movement and the back movement. A movement perpendicular to the shogging direction is likewise necessary between the back movement and the forth movement of the following loop formation operation.\nThe movement perpendicular to the shogging direction is often caused in that the guide bar is attached to a suspension shaft, which is rotated back and forth via a machine transmission. The shogging movement is controlled by a shogging drive, which moves the guide bar back and forth. The shogging drive is responsible among other things for the patterning, i.e., how large the lift of the guide bar is in each case in a loop formation operation is controlled with the aid of the shogging drive.\nA known method for controlling the shogging movement of the guide bar in a knitting machine is to use a pattern disk against which a contact roller bears. The contact roller is connected to a shogging slider. A ball pin is provided at the other end of the shogging slider, which ball pin is in engagement with a push rod, which in turn is in engagement with a ball pin on the guide bar. The push rod has ball sockets at both of its ends. The ball pins are thus hinge-mounted in the ball sockets. Thus, it is possible to also drive the guide bar through the pattern disk when it performs a swing-through motion perpendicular to the shogging direction. Through the two ball pin/ball socket bearings, the push rod compensates for the different distances that result between the front and the rear swing-through position on the one hand and the center position on the other hand.\nIn order to compensate for these changes caused by the swing-through movement of the effective length of the push rod, it has been proposed in DE 41 27 344 A1 to overlap a ground function with a compensation function. This compensation function must be taken into account when producing the pattern disk.\nOne disadvantage of this approach is still that alternately lateral forces act on the bar. This causes a negative alternating stress on the guide bar guidance and also generates frictional heat and wear in the ball connections. Another disadvantage is that the ball connections can be produced with the necessary precision only with relatively high expenditure. Even with a very precise production, temperature changes, which result, e.g., from the considerable frictional heat in the ball connections, cause division inaccuracies that lead to problems."}
{"text": "Hydrocracking is a versatile petroleum refining process, which enjoys widespread and growing use in the refining industry. Hydrocracking has the ability to process a wide range of difficult feedstocks into a variety of desirable products. Feedstocks, which may be treated by this process, include heavy naphthas, kerosenes, refractory catalytically cracked cycle stocks and high boiling virgin and coker gas oils. At high severities, hydrocracking can convert these materials to gasoline and lower boiling paraffins; lesser severities permit the higher boiling feedstocks to be converted into lighter distillates such as diesel fuels and aviation kerosenes.\nHydrocracking is conventionally carried out at moderate temperatures of 350° C. to 450° C. (650° F. to 850° F.) and at high pressures, over 7,000 kPa (1,000 psig) because the thermodynamics of the hydrocracking process become unfavorable at higher temperatures. In addition, high hydrogen pressures, usually at least 800 psig are usually required to prevent catalyst aging and so to maintain sufficient activity to enable the process to be operated with a fixed bed of catalyst for periods of one to two years without the need for regeneration.\nEnvironmental concerns, especially with sulfur oxides and nitrogen oxides emissions, have led petroleum refiners to depend more heavily than in the past on hydrodesulfurization and hydrocracking processes. Availability of by-product hydrogen from naphtha reforming no doubt has also cooperated to foster this dependence. Other factors too, have come into play to make hydroprocessing of increasing importance. Among these factors is that high quality crude oils for lube and fuels refineries are expected to progressively become more scarce. Also, refineries that include a fluid catalytic cracking (FCC) plant generate large volumes of dealkylated, aromatic refractory effluents, commonly known as FCC Cycle Oils. Decrease in demand for the fuel oil products into which these FCC Cycle Oils were previously incorporated has added to the practice of working them off by incorporation with a hydrocracker feedstock. The hydrocracking process, unlike catalytic cracking, is able to effectively upgrade these otherwise refractory materials.\nHydrocracking is an established petroleum refining process. The hydrocracking feedstock is invariably hydrotreated before being passed to the hydrocracker in order to remove sulfur and nitrogen compounds as well as metals and, in addition, to saturate olefins and to effect a partial saturation of aromatics. The sulfur, nitrogen and oxygen compounds may be removed as inorganic sulfur, nitrogen and water prior to hydrocracking although interstage separation may be omitted, as in the Unicracking-JHC process. Although the presence of large quantities of ammonia may result in a suppression of cracking activity in the subsequent hydrocracking step, this may be offset by an increase in the severity of the hydrocracking operation.\nIn the hydrotreater, a number of different hydrogenation reactions take place including olefin and aromatic ring saturation but the severity of the operation is limited so as to minimize cracking. The hydrotreated feed is then passed to the hydrocracker in which various cracking and hydrogenation reactions occur.\nIn the hydrocracker, the cracking reactions provide olefins for hydrogenation while hydrogenation in turn provides heat for cracking since the hydrogenation reactions are exothermic while the cracking reactions are endothermic; the reaction generally proceeds with generation of excessive heat because the amount of heat released by the exothermic hydrogenation reactions usually is much greater than the amount of heat consumed by the endothermic cracking reactions. This surplus of heat causes the reactor temperature to increase and accelerate the reaction rate, but control is provided by the use of hydrogen quench.\nConventional hydrocracking catalysts combine an acidic function and a hydrogenation function. The acidic function in the catalyst is provided by a porous solid carrier such as alumina, silica-alumina, or by a composite of a crystalline zeolite such as faujasite, Zeolite X, Zeolite Y or mordenite with an amorphous carrier such as silica-alumina. The use of a porous solid with a relatively large pore size in excess of 7 A is generally required because the bulky, polycyclic aromatic compounds which constitute a large portion of the typical feedstock require pore sizes of this magnitude in order to gain access to the internal pore structure of the catalyst where the bulk of the cracking reactions take place.\nThe hydrogenation function in the hydrocracking catalyst is provided by a transition metal or combination of metals. Noble metals of Group VIIIA of the Periodic Table, especially platinum or palladium may be used, but generally, base metals of Groups IVA, VIA and VIIIA are preferred because of their lower cost and relatively greater resistance to the effects of poisoning by contaminants. The preferred base metals for use as hydrogenation components are chromium, molybdenum, tungsten, cobalt and nickel; and, combinations of metals such as nickel-molybdenum, cobalt-molybdenum, cobalt-nickel, nickel-tungsten, cobalt-nickel-molybdenum and nickel-tungsten-titanium have been shown to be very effective and useful.\nThe porous support, which provides the acidic functionality in the catalyst, may comprise either an amorphous or a crystalline material or both. Amorphous materials have significant advantages for processing very high boiling feeds which contain significant quantities of bulky polycyclic materials (aromatics as well as polynaphthenes) since the amorphous materials usually possesses pores extending over a wide range of sizes and the larger pores, frequently in the size range of 100 to 400 Angstroms are large enough to provide entry of the bulky components of the feed into the interior structure of the material where the acid-catalyzed reactions may take place. Typical amorphous materials of this kind include alumina and silica-alumina and mixtures of the two, possibly modified with other inorganic oxides such as silica, magnesia or titania.\nCrystalline materials, especially the large pore size zeolites such as zeolites X and Y, have been found to be useful for a number of hydrocracking applications since they have the advantage, as compared to the amorphous materials, of possessing a greater degree of activity, which enables the hydrocracking to be carried out at lower temperatures at which the accompanying hydrogenation reactions are thermodynamically favored. In addition, the crystalline catalysts tend to be more stable in operation than the amorphous materials such as alumina. The crystalline materials may, however, not be suitable for all applications since even the largest pore sizes in these materials, typically about 7.4 Angstroms in the X and Y zeolites, are too small to permit access by various bulky species in the feed. For this reason, hydrocracking of residuals fractions and high boiling feeds has generally required an amorphous catalyst of rather lower activity.\nThe crystalline hydrocracking catalysts generally tend to produce significant quantities of gasoline boiling range materials (approximately 330° F.-, 165° C.-) materials as product. Since hydrocracked gasolines tend to be of relatively low octane and require further treatment as by reforming before the product can be blended into the refinery gasoline pool, hydrocracking is usually not an attractive route for the production of gasoline. On the other hand, it is favorable to the production of distillate fractions, especially jet fuels, heating oils and diesel fuels since the hydrocracking process reduces the heteroatom impurities characteristically present in these fractions to the low level desirable for these products.\nU.S. Pat. No. 8,372,772, issued Feb. 12, 2013, discloses a hydrocracking catalyst comprising zeolite crystallized as a layer on the surface of a porous alumina-containing matrix, the matrix being shaped by extrusion and arranged in a configuration to provide macropores in which the zeolite layer is provided on the walls of the macropores. The in-situ produced Y-fauajasite on matrix provides acidic functionality in the catalyst and is also used as a support for hydrogenating metals.\nU.S. patent applications (Pub. Nos. 2004/0220046 A1 and 2004/0235642 A1) teach preparation methods for structurally enhanced zeolite microsphere FCC catalysts where zeolite crystallizes on the walls of the pores of the matrix.\nAs is known in the industry, the hydrocracking process is a fixed-bed process with catalysts in shaped forms of extrudates, pellets or spheres. The particle size is much larger than FCC microspheres. While the above-mentioned U.S. Pat. No. 8,372,772 discloses in-situ crystallization of pre-shaped particles, such as extrudates, the patent does not otherwise provide a specific method which provides sufficient exchange of nutrients between the liquid phase and the solid phase to provide crystallization and allows the zeolite to form on the pore walls, as well as maintain the porosity and integrity of the extrudate."}
{"text": "Some missile systems require the use of vans which must accompany the missile in the field. Some of these vans carry material such as repair parts and test equipment which are necessary to the successful operation of the missile. Storage means must be provided in the vans for safely storing the material. The storage means typically is in the form of cabinets having a plurality of drawers therein which must be tightly secured against accidental opening while the vans travel over rough terrain with the missile as the missile is positioned in the field.\nCabinets are manufactured with built-in latches for securing the drawers in secured relation. However such latches are not capable of securing the drawers when the cabinets are subjected to vibrations such as encountered when the cabinets are mounted in vans and moved over rough terrain. In the past, the drawers have been additionally secured by various type of locking devices including barn door type of locking devices (hasps and staple) or the like. Usually such type of locking devices include a separate padlock or pin for attachment to each individual drawer for secured relation thereof to the cabinet. Each individual locking device must include a loosely carried pin or padlock.\nIt is an object of the present invention, therefore, to provide a locking device which is self storing.\nIt is a further object of the present invention to provide such a locking device which will maintain drawers in secured relation in cabinets while said cabinets are transported over rough terrain"}
{"text": "1. Field of the Invention\nThe present invention relates to a detecting device for detecting a transfer object, such as a document sheet or recording paper, in a facsimile machine, a printer, a copying machine or an image scanner for example.\n2. Description of the Related Art\nIn a facsimile machine for example, it is necessary to detect the movement or presence of a transfer object, such as a document sheet or a recording paper sheet, in a transfer path for controlling the operation of the facsimile machine. For this purpose, a detecting device is provided which comprises an actuator and a sensor.\nMore specifically, as shown in FIG. 8, a prior detecting device comprises an actuator 100 and a sensor 101. The actuator 100 includes a horizontal shaft 100a, a first arm 100b extending upwardly from the shaft 100a, and a second arm 100c extending downwardly from one end of the shaft 100a. The sensor 101 is mounted on a circuit board 102.\nIn the absence of a transfer object (not shown) in the transfer path, the first arm 100b of the actuator 100 projects into the transfer path. On the other hand, when a transfer object is present in the transfer path, the first arm 100b comes into engagement with the transfer object and is thereby pressed downwardly. In this way, the actuator 100 pivots up and down about the horizontal shaft 100a in response to the presence and absence of a transfer object.\nThe second arm 100c of the actuator 100 is bent to extend first obliquely and then substantially vertically. Since the second arm 100c is thus bent, the actual orientation of the second arm may be represented by an effective length line EL which passes through the center of the actuator shaft 100a and the lower end of the second arm 100c. In the absence of a transfer object in the transfer path, the angle .theta. between the effective length line EL of the second arm 100c and a vertical line VL passing through the center of the shaft 100a is less than 45.degree..\nThe lower end of the second arm 100c is provided with a reflector 103 which faces the sensor 101 when the second arm 100c is pivoted down due to the absence of a transfer object. Further, the reflector 103 has a rounded stopper projection 104 which comes into resting contact with the circuit board 102, thereby maintaining a predetermined clearance H.sub.1 between the reflector 103 and the circuit board 102.\nThe sensor 101 is a reflection type sensor which has a light emitting portion and a light receiving portion. When the actuator 100 is pivoted to bring the stopper projection 104 into contact with the circuit board 102, the light emitted from the light emitting portion of the sensor 101 is reflected on the reflector 103 for incidence into the light receiving portion of the sensor 101, so that the sensor 101 notifies the absence of a transfer object. On the other hand, when the second arm 100c of the actuator 100 is pivoted up due to the presence of a transfer object, the reflector 103 becomes far from the sensor 101 to come completely out of the light reflecting position (see the phantom line in the center of FIG. 8), so that the sensor 101 notifies the presence of a transfer object.\nNormally, the circuit board 102 is mounted at a standard distance H.sub.0 from the actuator shaft 100a to make the reflector 103 substantially parallel to the upper surface of the sensor 101 in facing relation thereto (see the center representation in FIG. 8) when the stopper projection 104 of the reflector 104 comes into contact with the circuit board 102. However, it is possible, due to some production or assembly error, that the circuit board 102 may positionally deviate from the standard distance position H.sub.0 relative to the actuator shaft 100a. If the circuit board 102 is erroneously positioned at a distance (H.sub.0 -.DELTA.H) or (H.sub.0 .DELTA.H) from the actuator shaft 100a (see the left and right representations in FIG. 8), the orientation of the reflector 103 becomes improper relative to the upper surface of the sensor 101. As a result, the sensor 101 may fail to provide a reliable detection of a transfer object. Such a disadvantage will become particularly remarkable as the angle .theta. between the effective length line EL of the actuator second arm 100c and the vertical line VL is small."}
{"text": "The present invention relates to photopolymer materials sensitive to infrared, near infrared, red and green light radiation for initiating polymerization and to applications of such photopolymer, like holographic polymer dispersed liquid crystal (H-PDLC) or reversible dye doped photopolymer (RDDP) materials, for making optical devices. The invention relates to holographic polymer dispersed liquid crystal and reversible dye materials having improved electrical and optical switching properties.\nInfrared (IR) diode sources are largely used in integrated photonic circuits and the list of their applications grows very rapidly. The management of the radiation of these sources requires the fabrication of optical elements such as lenses, interconnects, modulators, etc. Holographic diffractive elements are emerging as very promising for these applications. These include wavelength selective holographic interconnects (1 to N, or M to N), couplers, lenses, mirrors. Today, a new class of Holographic Photopolymer Dispersed Liquid Crystal Materials (H-PDLCs) is considered to be one of the most viable technologies for the development of reflective color displays, switchable holographic optical elements (such as Bragg gratings for Wavelength Division Multiplexing (WDM) devices), switchable-focus lenses, etc. See for example Crawford et al., xe2x80x9cReflective color LCDs based on H-PDLC and PSCT technologiesxe2x80x9d, J Soc. Information Display, 1996, 5(1); Domash et al., xe2x80x9cElectronically switchable waveguide Bragg gratings for WDM routingxe2x80x9d, 1997 Digest of the IEEE. /LEOS Summer Topical Meetings: Vertical-Cavity Lasers/Technologies for a Global Information Infrastructure/WDM Components Technology; and Domash et al., xe2x80x9cSwitchable-focus lenses in holographic polymer dispersed liquid crystalxe2x80x9d, Proceedings of the SPIExe2x80x94The International Society for Optical Engineering, vol. 2689, (Diffractive and Holographic Optics Technology III, San Jose, Calif., USA, 1-2 Feb. 1996.) SPIE-Int. Soc. Opt. Eng, 1996. p. 188-94.\nCommercially available holographic materials are generally sensitive in the UV/visible part of the spectrum only. Thus, the actual fabrication of above-mentioned elements requires an initial recording step, where a UV-Visible laser source is used, and then, a further adaptation or adjustment of the obtained element for utilization with near IR wavelengths (e.g. 800-850 nm or 1300-1500 nm), which are used in local or long distance communication systems. Due to strong astigmatism and divergence of the used diode lasers, this work is difficult and has poor efficiency. It would therefore be highly desirable to provide in situ recording of holographic diffractive elements with lasers, which are already integrated in the given photonic circuit, thus providing self-alignment of the photonic circuits. Thus, there is a need to extend the sensitivity of these materials up to communication wavelengths. Namely, for in situ holographic recording of optical components with diode lasers operating in the 800-855 nm, among which there are Vertical Cavity Surface Emitting Lasers (VCSEL), it is important to have H-PDLC material with suitable holographic characteristics, which include high sensitivity and diffraction efficiency of the recording, low scattering and noise level, and low switching voltage.\nUp to now, the sensitivity of the existing H-PDLC materials has been extended up to 790 nm (Natarajan et al., Report on Photonics West 98, San Jose, January 1998). Photopolymerizable materials have been recently proposed with sensitivity in the 800-850 nm area, for example in co-pending application Ser. No. 09/503,207, filed on Feb. 14, 2000. See also EP 0 223 587; EP 0 387 087; EP 0 389 067; U.S. Pat. No. 4,343,891; Chatterjee et al., J. Am. Chem. Soc., 1990, 112, 6329; Schuster et l, Photochem. Photobiol. A: Chem., 1992, 65, 191; Chatterjee et al., J. Am. Chem. Soc., 1988, 110, 2326; Cooper et al., J. Am. Chem. Soc., 1963, 85, 1590; and Noiret et al., Pure and Applied Optics, 1994, 3(1), 55-71. Imaging applications with low resolution (such as printing plates) were also successfully explored. Some of the materials were the subject of the study for holographic gratings recording and only very low level of performance was achieved (7% of diffraction efficiency at sensitivity of about 300-500 mJ/cm2), as reported in the mentioned Noiret et al. article supra. This low level of performance makes such polymers impractical for commercial application. All the above materials still suffer from limitations such as spatial resolution, diffraction efficiency etc.\nIn accordance with the present invention, there is now provided a photopolymerizable formulation sensitive to light in a green to infrared region of the optical spectrum comprising:\na) a photopolymerizable monomer, oligomer or mixtures thereof;\nb) a photoinitiator sensitive to light in the green to infrared region;\nc) an additive for increasing the refractive index and decreasing the viscosity of the formulation; and\nd) an optional filler having optical properties selected to contrast with optical properties of a polymer resulting from photopolymerization of the monomer, oligomer or mixtures thereof.\nPreferably, the formulation comprises filler, which is a liquid crystal or a reversible dye, or combinations thereof, or any other filler materials such as mesogens having polar or functional groups. The liquid crystal preferably has a polar group attached to ends of the molecule chain elements, whereby droplets of LC matter are more efficiently formed in the polymer dispersed in the liquid crystal material. Larger droplets allow for more efficient optical state switching.\nIn a second aspect of the invention, there is provided a process for making an optical device comprising the steps of:\npreparing a photopolymerizable formulation sensitive to light from green to infrared region of the optical spectrum, the formulation comprising\na) a photopolymerizable monomer, oligomer or mixtures thereof;\nb) a photoinitiator sensitive to light in the region\nc) an additive for increasing the refractive index and decreasing the viscosity of the formulation; and\nd) an optional filler having optical properties selected to contrast with optical properties of a polymer resulting from photopolymerization of the monomer, oligomer or mixtures thereof;\napplying a layer of the formulation on an optical element;\nfirst exposing the optical element comprising the formulation to a light source emitting light in the region, to polymerize the formulation and produce a recording pattern on the optical element; and\noptionally second exposing the optical element to the light to polymerize any remaining portion of the formulation not polymerized during first exposing, thereby producing said optical device.\nFinally, in a third aspect of the present invention, there is provided a process of making an optical connection between at least two waveguides or light guides of an optical element by photopolymerization, the process comprising the steps of:\npreparing a photopolymerizable formulation sensitive to light from green to infrared region of the optical spectrum, the formulation comprising:\na) a photopolymerizable monomer, oligomer or mixtures thereof;\nb) a photoinitiator sensitive to light in the green to infrared region;\nc) an additive to increase the refractive index and decrease the viscosity of the formulation; and\nd) an optional filler;\napplying a sufficient amount of the formulation between the waveguides or light guides of the optical element to be connected; and\nexposing the optical element to a light source emitting light in the region to polymerize the formulation, thereby making an optical connection between the at least two waveguides or light guides.\nPreferably, the light is transmitted during the exposing through at least one of the at least two waveguides or light guides. While a filler is also optional in the case of optically coupling waveguides or light guides, a filler having the desired optical properties is preferred to reduce the amount of polymer used and to impart desirable physical properties to the formulation."}
{"text": "A CD-ROM appendix is included in this disclosure. Specifically, Appendix B is a plurality of tables utilized by the computer source code listing. The CD-ROM is submitted at the same time as this preliminary amendment, and is hereby incorporated by reference. The only file on the CD-ROM is entitled, xe2x80x9c10932-43 CD-ROM Appendix.xe2x80x9d The file size is 790 KB and the file was created on Nov. 27, 2001. The machine format is IBM-PC and the operating system used to create the file is MS-Windows.\n1. Technical Field\nThe present invention relates generally to speech encoding and decoding in voice communication systems; and, more particularly, it relates to various techniques used with code-excited linear prediction coding to obtain high quality speech reproduction through a limited bit rate communication channel.\n2. Related Art\nSignal modeling and parameter estimation play significant roles in communicating voice information with limited bandwidth constraints. To model basic speech sounds, speech signals are sampled as a discrete waveform to be digitally processed. In one type of signal coding technique called LPC (linear predictive coding), the signal value at any particular time index is modeled as a linear function of previous values. A subsequent signal is thus linearly predictable according to an earlier value. As a result, efficient signal representations can be determined by estimating and applying certain prediction parameters to represent the signal.\nApplying LPC techniques, a conventional source encoder operates on speech signals to extract modeling and parameter information for communication to a conventional source decoder via a communication channel. Once received, the decoder attempts to reconstruct a counterpart signal for playback that sounds to a human ear like the original speech.\nA certain amount of communication channel bandwidth is required to communicate the modeling and parameter information to the decoder. In embodiments, for example where the channel bandwidth is shared and real-time reconstruction is necessary, a reduction in the required bandwidth proves beneficial. However, using conventional modeling techniques, the quality requirements in the reproduced speech limit the reduction of such bandwidth below certain levels.\nSpeech encoding becomes increasingly difficult as transmission bit rates decrease. Particularly for noise encoding, perceptual quality diminishes significantly at lower bit rates. Straightforward code-excited linear prediction (CELP) is used in many speech codecs, and it can be very effective method of encoding speech at relatively high transmission rates. However, even this method may fail to provide perceptually accurate signal reproduction at lower bit rates. One such reason is that the pulse like excitation for noise signals becomes more sparse at these lower bit rates as less bits are available for coding and transmission, thereby resulting in annoying distortion of the noise signal upon reproduction.\nMany communication systems operate at bit rates that vary with any number of factors including total traffic on the communication system. For such variable rate communication systems, the inability to detect low bit rates and to handle the coding of noise at those lower bit rates in an effective manner often can result in perceptually inaccurate reproduction of the speech signal. This inaccurate reproduction could be avoided if a more effective method for encoding noise at those low bit rates were identified.\nAdditionally, the inability to determine the optimal encoding mode for a given noise signal at a given bit rate also results in an inefficient use of encoding resources. For a given speech signal having a particular noise component, the ability to selectively apply an optimal coding scheme at a given bit rate would provide more efficient use of an encoder processing circuit. Moreover, the ability to select the optimal encoding mode for type of noise signal would further maximize the available encoding resources while providing a more perceptually accurate reproduction of the noise signal.\nA random codebook is implemented utilizing overlap in order to reduce storage space. This arrangement necessitates reference to a table or other index that lists the energies for each codebook vector. Accordingly, the table or other index, and the respective energy values, must be stored, thereby adding computational and storage complexity to such a system.\nThe present invention re-uses each table codevector entry in a random table with xe2x80x9cLxe2x80x9d codevectors, each of dimension xe2x80x9cN.xe2x80x9d That is, for example, an exemplary codebook contains codevectors V0, V1, . . . , VL, with each codevector Vx being of dimension N and having elements C0, C1, . . . , CN-1, CN. Each codevector of dimension N is normalized to an energy value of unity, thereby reducing computational complexity to a minimum.\nEach codebook entry essentially acts as a circular buffer whereby N different random codebook vectors are generated by specifying a starting point at each different element in a given codevector. In one embodiment, each of the different N codevectors then has unity energy.\nThe dimension of each table entry is identical to the dimension of the required random codevector and every element in a particular table entry will be in any codevector derived from this table entry. This arrangement dramatically reduces the necessary storage capacity of a given system, while maintaining minimal computational complexity.\nOther aspects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "A slender structure is very sensitive to dynamic loading. A high building, long span bridge, skyscraper and other similar structures may lose their function or in some serious cases may be destroyed under dynamic loading. In general, the kind of damage induced by the dynamic loading can be divided into two kinds, namely functional and safety loss.\nIn order to solve this kind of problem, it is conventional to increase the strength and ductility of load bearing structures. This is the usual practice in earthquake-resistance design for civil works. Another method is to employ some measure to reduce the response of the structure to such vibration.\nThe methods of reducing vibration of the structure generally can be divided into two classes, namely passive control and active control. The so-called passive control is usually carried out by increasing or decreasing a number of springs in the structure to change the stiffness of the structure, so that its natural frequency can be adjusted to depart from the frequency of dynamic loading.\nAlternatively, a mass-spring damper may be installed on the structure to concentrate vibration energy in the mass-spring damper, so that the goal of reducing the vibration of main structure can be reached.\nAnother method is to arrange weak spots at some specific positions on the structure. In this way, strong vibrations will force these weak spots to produce plastic strains, by which energy consumption of the total structure can be increased. A common point of the passive control is that the dynamic characteristics of the structure are fixed after completing the structure, and the stress status of struts in the structure is completely determined by the time history of the loading. It is not necessary or even possible to take initiative action to change the characteristics. Therefore the designer must grasp in advance the complete characteristics of dynamic loading, such as magnitude and spectrum, so that an effect of the vibration reduction by these methods can be assured.\nThe so-called active control generally means installing a mechanism, such as an oil or air pressure mechanism, which can give additive force to the structure or change the status of inner forces in a part of the struts in the structure. Though theoretically the active control method can achieve a very good effect, it has a problem of stability in actual application, and also a large amount of energy has to be supplied. Furthermore, the mechanism itself is unreliable. Its range of applications, and especially applications for reducing the vibration of civil engineering and construction under dynamic loading, are much less than passive control cases.\nAs mentioned above passive control has the fault that its effect is poor, and active control has the faults that it is unreliable and requires a large amount of energy."}
{"text": "A variety of automated or semi-automated chemical analyzers are known which utilize cuvettes for the chemical testing of samples placed therein. Generally, a predetermined amount of liquid sample, such as biological fluid is placed in the cuvette which is then transported through the instrument. As the cuvette is transported, the instrument dispenses a quantity of reagent into the sample and monitors the resulting chemical reaction. Such monitoring is generally accomplished through use of an optical means which views the fluid sample through optically transparent portions of the cuvette.\nIn order to simplify the loading of the cuvettes into the instrument and facilitate their handling by the instrument once so loaded, proposals have been made to provide the cuvettes in the form of a continuous integral strip. The individual cuvettes of the strip are designed to be relatively rigid but the strip itself is provided with sufficient flexibility to ease its transport through the instrument.\nFurthermore, by making the cuvettes in a continuous strip form, they can be manufactured relatively inexpensively from suitable plastic material, thereby permitting their disposal after use. This is an important feature since it avoids the requirement for washing the cuvettes after use and avoids any possibility of cross-contamination of fluid samples which could cause erroneous test results. A proposed cuvette system designed to meet these requirements is proposed in U.S. Pat. No. 4,263,256.\nIn commonly owned copending U.S. patent application Ser. No. 284,842 filed July 20, 1981 and entitles \"Cuvette System For Automated Chemical Analyzer\", now abandoned, the disclosure of which is hereby incorporated by reference in its entirety herein, there is described a cuvette belt which comprises a matching pair of formed elongated plastic strips which are joined together along corresponding faces thereof to form an integral belt. A series of regularly spaced cavities or chamber halves is formed transversely in each of the corresponding strip faces which define open-topped cuvette receptacles when the belt halves are joined. As described, the cuvette belt is made by forming strip plastic material with a series of regularly spaced transverse (laterally extending) elongated pockets so as to define two integral side-by-side belt halves. The formed strip is then divided longitudinally to separate the belt halves and the belt halves brought into register and joined together to form a completed cuvette belt.\nUsing such manufacturing techniques, cuvettes may be obtained with superior optical characteristics and dimensional accuracy previously obtainable only through use of injection molding techniques. This is important when optically analyzing samples in the cuvettes for which a precisely defined optical path through the cuvette is required. It is pointed out that it is possible by utilizing cold forming techniques to avoid optical degradation of the material due to heat. Further an optical portion of the material may be restrained by clamping or other forming techniques during the pocket forming operation to avoid stretching or other deformation of portions of the pocket which form the sides of the cuvette. In this manner essentially all stretching of the material during forming is limited to the side walls of the pockets and optical portions thereof are maintained stress-free and with a uniform thickness.\nReference is also made to commonly owned copending U.S. patent application Ser. No. 746,231 filed June 18, 1985, and entitled \"Cuvette Belt Manufacture and Process\", the disclosure of which application is hereby incorporated by reference in its entirety herein. In that application is disclosed an alternative method for making cuvette belts where they comprise a matching pair of plastic strips, each of which is formed with chamber halves, and which are joined together to form an integral cuvette belt with the chamber halves aligned to form the cuvettes.\nAccording to that method, two strips of plastic material are identically formed with a series of regularly spaced transverse elongated pockets so as to define two integral side-by-side belt halves. The two formed strips are brought into register and joined together to form a composite strip defining two integral mirror image cuvette belts joined by their cuvette mouth ends. The composite strip is then divided longitudinally to separate the cuvette belts.\nThe present invention is particularly concerned with techniques for joining the formed plastic strips by heat sealing during the manufacture of such cuvettes while retaining the aforesaid optical characteristics."}
{"text": "1. Field of the Invention\nThe present invention relates to improvement in optical gas analyzers and more particularly, to refinements in non-dispersive infrared analyzers in order to obtain adequate performance for medical and other precision applications.\n2. Description of Prior Art\nThe inventor of the present invention has contributed to the prior art devices in the area of infrared gas analyzers and this prior art as exemplified in U.S. Pat. No. 3,932,754 substantially sets forth one aspect of devices commercially utilized to date for non-medical applications. The inventor in setting forth to perfect instrumentation for medical applications found that the prior art teachings contained in the above referenced patent and others known in the art, lacked the degree of reliability and sophistication necessary for medical applications.\nTowards this end, when dealing with instrumentation in life support systems in which the output of information is of a critical nature the degree of calibration and other performance standards, as hereinafter explained in detail, was not available with the prior art equipment. Accordingly, although the present invention will be particularly discussed with respect to the benefits it provides for medical application, these same advantages also lend themselves to industrial and other applications as well."}
{"text": "Consumer demand for Internet-based services has led to the widespread deployment of customer premise equipment (CPE) to provide various connectivity functions. The growth of home networking technologies has further expanded the range of devices installed at customer premises. The increased amount of voice, video and data traffic is driving the market for premise-based services. This growth is only expected to increase over time and the resulting costs of maintaining a complex assortment of CPE is also expected to increase. The cost of configuring new features and services on hardware network appliances is particularly burdensome because of the cost of sending technicians to customer sites.\nBased on the foregoing, there is a need for an efficient approach to the remote management and support of network appliances installed at customer premises."}
{"text": "1. Field of the Invention\nThe present invention relates to a hose-in-hose or double walled hose, and, more particularly, to a hose-in-hose secondary coupler.\n2. Description of the Related Art\nIt is known that hose-in-hoses have been used to transport liquid waste, such as nuclear waste, to a place away from the source.\nA problem with current hose-in-hose configurations is that once the outer hose is disconnected via a coupler, the separated ends of the outer hoses must be physically pulled apart to allow access to the inner hose. Although on shorter lengths of hose this may not present much of a problem, as the length of the hose increases the difficulty in providing adequate force to pull apart the disconnected ends also increases. The ability to adequately separate the disconnected ends of the outer hose to a significant distance thereby allowing access to the inner hose relies on the strength of the end user. Furthermore, longer hoses may not be able to be separated at the coupler due to the weight of the individual hose ends.\nWhat is needed is a device which allows access to the inner hose without the need to substantially move the two disconnected ends of the outer hose."}
{"text": "Tn U.S. Pat. No 3,612,793 issued on Oct. 12, 1971, to John O. Roeser and entitled \"Electrical Switch Components And Switches Formed Thereby,\" there is shown and described a unitary combination switchblade and contact means For compact switch mechanisms. As can be seen from FIGS. 8, 9 and 10 of the '793 patent, the combination switchblade and contact means 16 comprises a pair of identical members 109a (109b) each being formed of a first end portion 110 a mid-portion 116 and a curvilinear second end portion 118. The first end portion 110 consists of a pair of spaced apart legs 112 and 114. The legs are spaced apart by a dimension D so as to receive therebetween a tension spring 120, as illustrated in FIG. 5. A combination strut and attachment means 124 is shear formed in the curvilinear second end portion 118 and includes a reverse curve portion 131 for mating with the end of the tension spring 120 (FIG. 5).\nThe curvilinear second end portion 118 further includes an upper contact portion 128 and a lower contact portion 130 which is separated from the upper contact portion 128 by means of the reverse curve portion 131. The legs 112 and 114 are provided with sharpened edges 132 for pivoting within opposed V-shaped notch portions 106 formed within an elongated metal member 74 of the actuating means 114. The lower contact portion 130 is located within the same plane as the first portion 110 and the mid-portion 116 However, the upper contact portion 128 is off-set from that same plane. The lower and 15 upper contact portions 128 and 130 define movable contact members which are adapted to be snapped between opposed fixed contact portions upon upper terminal members 134 and lower terminal members 136, as illustrated in FIG. 5.\nThe upper portions of the upper terminal members 134 are formed with an off-set transverse portion 116 having a depending upper fixed contact 148 (FIGS. 14 and 15). The upper fixed contacts 148 serve to make electrical contact engagement with the upper contact portions 128 of the movable contact members. The lower terminal members 136 are provided with portions 156 which serve as lower fixed contacts for electrical contact engagement with lower contact portions 130 of the movable contact members.\nThe combination switchblade and contact means 16 is assembled with other components so as to form a compact precision, snap-action push-button switch 10 and is illustrated in FIGS. 2, 5 and 6. In operation, when the button 72 of the actuator means 14 is depressed so as to move downwardly, this causes the elongated metal member 74 to also more downwardly against the bias of the return springs 76 and 78. Initially, the combination 15 switchblade and contact means is restrained by means of the lower fixed contact portions 156 so that the lower contact portions 130 of the movable contact members remain in contact therewith.\nThis contact continues until the pivot points 106 pass the center line of the tension spring 120. At that time, the movable contact members will snap overcenter and the upper contact portions 130 thereof will become engaged with the upper fixed contacts 148. This engagement continues until such time when the pressure upon the top of the button 72 is released, thereby causing the return springs 76 and 78 to return the metal member 74 and button 72 to their rest positions shown in FIGS. 1 and 2. As the metal member 74 returns past the center line of the tension spring 120, the movable contact members will snap overcenter again in the reverse direction.\nThe present invention represents an improvement over the combination switchblade and contact means described above in connection with the '793 patent. The upper contact portion 128 and the lower contact portion 130 upon the second end portion 118 of the switchblade and contact means 16 in the aforenoted patent experience a high degree of stress during the switching operations which tends to shorten the service lives of such components. It would therefore be desirable to provide a novel and improved switch butterfly assembly for use in precision snap action switches which has a unique configuration so as to reduce the stresses upon the movable contact portions, thereby prolonging the service lives thereof."}
{"text": "Melanomas are aggressive, frequently metastatic tumors derived from either melanocytes or melanocyte related nevus cells (“Cellular and Molecular Immunology” (1991) (eds.) Abbas A. K., Lechtman, A. H., Pober, J. S.; W. B. Saunders Company, Philadelphia: pages 340-341). Melanomas make up approximately three percent of all skin cancers and the worldwide increase in melanoma is unsurpassed by any other neoplasm with the exception of lung cancer in women (“Cellular and Molecular Immunology” (1991) (eds.) Abbas, A. K., Lechtiman, A. H., Pober, J. S.; W. B. Saunders Company Philadelphia pages: 340-342; Kirkwood and Agarwala (1993) Principles and Practice of Oncology 7:1-16). Even when melanoma is apparently localized to the skin, up to 30% of the patients will develop systemic metastasis and the majority will die (Kirkwood and Agarwala (1993) Principles and Practice of Oncology 7:1-16). Classic modalities of treating melanoma include surgery, radiation and chemotherapy. In the past decade, immunotherapy and other molecular methods have emerged as new and promising methods for treating melanoma.\nStrong evidence that an immune response to cancer exists in humans is provided by the existence of lymphocytes within melanoma deposits. These lymphocytes, when isolated, are capable of recognizing specific tumor antigens on autologous and allogeneic melanomas in a major histocompatibility complex (MHC)-restricted fashion (Itoh et al. (1986), Cancer Res. 46: 3011-3017; Muul et al. (1987), J. Immunol. 138:989-995); Topalian et al. (1989) J. Immunol. 142: 3714-3725; Darrow et al. (1989) J. Immunol. 142: 3329-3335; Hom et al. (1991) J. Immunother. 10:153-164; Kawakami et al. (1992) J. Immunol. 148: 638-643; Hom et al. (1993) J. Immunother. 13:18-30; O'Neil et al. (1993) J. Immunol. 151: 1410-1418). Tumor infiltrating lymphocytes (TIL) from patients with metastatic melanoma recognize shared antigens including melanocyte-melanoma lineage specific tissue antigens in vitro (Kawakami et al. (1993) J. Immunother. 14: 88-93; Anichini et al. (1993) J. Exp. Med. 177: 989-998). The fact that many melanoma patients mount cellular and humoral responses against these tumors and that melanomas express both MHC antigens and tumor associated antigens (TAA) suggests that identification and characterization of additional melanoma antigens will be important for immunotherapy of patients with melanoma. However, there remains a need for new modalities for the treatment of melanoma and other cancers."}
{"text": "This invention relates to electronic circuitry and, more particularly, to a charge pump for producing a negative substrate bias in a complementary metal oxide semiconductor (CMOS) integrated circuit.\nMOS transistors are commonly used in electronic circuits such as dynamic random access memories (DRAMS). In an NMOS transistor, an N-type source region is separated from an N-type drain region by a P-type channel region. All three regions are formed in a P-type semiconductor substrate. By applying a positive voltage to a gate electrode disposed above the channel region, electrons gather in the channel region between the source region and the drain region to allow current to flow from the drain region to the source region. PMOS transistors have the same structure except the conductivity types of the various regions are reversed and a negative gate voltage is required to allow current to flow from the source region to the drain region.\nIt has been found that NMOS transistors operate better when the P-type substrate of the NMOS (or of the NMOS transistors in a CMOS circuit) is driven negative with respect to circuit ground, in other words there is a negative substrate bias. Such a negative substrate bias provides a number of advantages in terms of the overall circuit performance. More specifically, a negative substrate bias decreases the NMOS transistor source and drain capacitance, decreases the likelihood of latchup, decreases PN diode injection when a node is driven below ground, and decreases the effective body effect, all of which are desirable in CMOS circuits.\nTypically a charge pump circuit is used to create the negative substrate bias. Once a negative substrate bias is achieved, however, it does not last forever. For example, when an NMOS transistor is conductive with a relatively high drain to source voltage, some of the electrons traveling from the source region to the drain region collide with atoms in the channel region with enough energy to cause electron/hole pairs to form. The positive gate voltage attracts the generated electrons to the surface of the channel while the positive drain voltage attracts them to the drain where they simply add to the normal flow of electrons from source to drain. The positively charged holes, by contrast, are repelled by the positively charged gate away from the channel region into the substrate. The substrate current created by the excess holes makes the substrate more positively charged, thus counteracting the negative substrate bias. In DRAMS, a substantial amount of substrate current is generated whenever the memory is read or written, since many transistors are switched on and off at that time. This component of substrate current may be orders of magnitude above the background (i.e., standby) leakage current of all the reverse biased P--N diodes throughout the circuit. Therefore, the charge pump must remove low substrate current during standby and high substrate current during high activity to maintain the negative substrate bias.\nFIG. 1 is a conceptual schematic diagram of a charge pump 2 which includes a first switch 4 coupled between a positive power supply voltage (V.sub.cc) and, a first terminal 6 of a capacitance C1. A second switch 8 is coupled between a ground potential (V.sub.ss) and a second terminal 10 of capacitance C1. A third switch 12 is coupled between (V.sub.ss) and terminal 6 of capacitance C1, and a fourth switch 14 is coupled between the substrate (represented by the voltage (V.sub.bb)) and terminal 10 of capacitance C1. In operation, switches 4 and 8 are both closed (made conductive) for charging capacitance C1 to a voltage equal to the difference between (V.sub.cc) and (V.sub.ss). In FIG. 1, (V.sub.cc)=+5 volts and (V.sub.ss)=0 volts, so capacitance C1 charges with node 6 five volts more positive than node 10. Thereafter, switches 4 and 8 are opened and switches 12 and 14 are both closed. Since the positive terminal 6 of capacitance C1 is now coupled to a ground potential, the negative terminal 10 of capacitance C1 tries to drive V.sub.bb to negative 5 volts through switch 14. Thereafter, switches 12 and 14 are opened, and the sequence repeats itself. An oscillator (not shown) typically controls the repetitive switching sequence, and a detector (not shown) monitors the substrate voltage and controls the pumping operation to maintain the substrate at the proper negative voltage level.\nAs discussed in more detail below, known charge pumps consume a substantial amount of power (often 1 milliwatt or more even when no further pumping is required), often work against themselves by adding positive substrate current as they operate, and generally operate inefficiently."}
{"text": "This invention relates in general to die-casting machines and in particular to high-speed die-casting machines for the production of precision miniature castings where the weight of the produced casting generally does not exceed two ounces and the face of the die generally does not exceed approximately ten square inches. While these size limitations provide a suitable focus, many of the features of the present invention are applicable to much larger machines.\nHeretofore, machines for the die casting of precision miniature pieces have been generally of the Dynacast type, a die-casting machine developed in the mid-1960's and the subject of several United States patents, including U.S. Pat. No. 3,364,981. This machine, at the time of its development, drastically changed the manufacture of miniature die-casting pieces which prior to the time of the U.S. Pat. No. 3,364,981 were manufactured in large conventional die-casting machines using multiple-cavity dies. However, the state of the art has remained practically unchanged during the last 20 years. At the expiration of the above-mentioned patent, the machine has been duplicated by several manufacturers with more or less the same features but without eliminating the various shortcomings.\nOne shortcoming of the Dynacast machine is the fact that the machine can use only two vertical cores and only occasionally a frontal core. The zinc injection system of the Dynacast machine is made by a gooseneck partially submerged into the liquid zinc of the melting pot. The gooseneck has a swiveling movement in order to bring the nozzle in contact with the die face and to move the nozzle away from the die face to break the runner from the liquid metal. Besides the mechanical elements needed for this movement (which adds complexity to the machine), this system reduces the speed at which the machine can operate. The position of the gooseneck precludes the use of a fourth core.\nAnother shortcoming of the Dynacast machine includes the way the level of the zinc in the neck of the gooseneck is controlled. This system leaves a certain amount of air ahead of the zinc in the neck of the gooseneck, which amount of air is variable and depends upon several parameters which are not easily controlled, such as the temperature of dies, size of the shots, speed of the machine and temperature of the nozzle. Another drawback of the Dynacast machine is the way the toggle mechanism operates in order to close the die halves and to create the force between them. This force generated between the dies varies, depending on the wear of the mechanical components of the toggles, temperature of the dies, temperature of the plate of the machine, and so forth.\nA further drawback of the Dynacast machine is the fact that it cannot produce parts with what is known in the industry as a \"hardware\" finish; in other words, parts that do not need to be plated. The reason for such inability involves the way the zinc is injected into the die, the gating (which has to be small because of the subsequent separation from the part with the tumbling operation) and the system of heat transfer on the die. The parts from the Dynacast machine are ejected with the gates attached to them, and this requires a subsequent operation (tumbling) in order to separate the parts from the gate, with all of the negative attributes of such operation including surface pitting, dimensional alteration, and breaking of sharp corners.\nThe present invention provides a machine which is designed to overcome the above shortcomings and disadavantages. In order to achieve this object, the disclosed machine produces castings which may require up to four cores along the plane of the dies (substantially two vertical and two horizontal) and two cores in a plane perpendicular to the surface of the die with a total of six cores. The machine of the present invention produces castings with much less porosity than any other machine used today and with surface finish of the type as previously mentioned which is known in the industry as a hardware finish. The cast pieces can be trimmed automatically in the machine so that the pieces are separated from the gate, the gates being discharged in a separate container for a subsequent remelt.\nAnother reference of possible interest to the heat transfer aspects of the present invention is U.S. Pat. No. 4,248,289 which issued to Perella et al. This patent discloses a means to heat the heat-exchange medium, water, by electric immersion heaters to preheat the die aove the normal boiling point of water by controlling the internal pressure of the die-cooling cavity. Die cooling is accomplished by the heat of evaporation at the temperature and pressure required for best casting quality. The wall thickness of the die cavity to water cavity is sized according to the heat-transfer rate or the die-cavity material. As steam is produced within the cavity, make-up water is introduced to maintain a wetted surface.\nIn die-casting practice, the \"coring\" of a cast part encompasses many forms. The core may be in the form of an entry hole, round, square or otherwise, a hollow space, a countersunk or extended boss, an embossed face, or any shape that is required that is not formed by the two-part mold cavity. Some core requirements are parallel to the line of separation of the mold halves but are independent in movement. Other cores not in the line of mold separation movement must be positioned and supported either by devices integral with the die blocks or as a function of the die-casting machine. Those arrangements which are integral with the molding dies add substantially to the cost of the dies and are not practical for miniature casting work. Since the hydraulic pressure area of some cores can be equal to that of the die cavity, it is necessary to have equal support to that of the die cavity to lock the core in a fixed position. This may be accomplished by toggle movement to rapidly position the core in the die and then lock the position as securely as the main die-closing action."}
{"text": "The present invention relates to a tapping fitting for application to a pipe in order to direct fluid conducted therethrough to an ancillary device, such as a branch conduit or an outlet tap. The invention is particularly useful in electrofusion fittings for plastic pipes wherein the fitting is bonded to a plastic pipe by electrofusion, and is therefore described below with respect to this application, although it will be appreciated that the invention could also be used in other applications.\nTapping fittings of the foregoing type are known. They generally include a rotatable blade assembly formed with external threads threadedly received within the housing and including a cylindrical cutting blade such that rotation of the blade assembly moves it axially from an initial, retracted position within the housing, to a projected position, to cut a cylindrical slug from the plastic pipe to which the fitting is attached. Such tapping fittings are applied to the plastic pipe while the blade assembly is in its initial, retracted position, i.e., within the housing of the fitting. After the fitting has been so applied, the blade assembly is rotated from externally of the fitting to its projected position to cut a cylindrical slug from the plastic pipe, and then it is rotated in the opposite direction back to its initial, retracted position to remove the slug and retain it within the housing main section of the fitting. The fitting thus establishes communication between the plastic pipe and the ancillary device, and generally remains in place for the life of the installation.\nSuch tapping fittings have come into extensive use for tapping branch conduits to high-pressure liquid or gas supply pipes of a plastic material. The housing is generally made of polyethylene or of a similar plastic material as the pipe. It has been found, however, that the plastic material used for the housing tends to creep over a long period Of time. Transverse creep (i.e., transversely to the axial movement of the blade assembly) may result in the blade assembly jumping the threads; whereas axial creep may result in the failure of the valve to completely close."}
{"text": "Ovarian cancer is the eighth most common type of cancer among women. The American Cancer Society estimates that about 22,220 new cases of ovarian cancer will be diagnosed in the United States during 2005. Ovarian cancer accounts for about 3% of all cancers in women. Because many ovarian cancers cannot be detected early in their development, they account for a disproportionate number of fatal cancers, being responsible for almost half the deaths from cancer of the female genital tract; more deaths than any other reproductive organ cancer. Older women are at higher risk. More than half of the deaths from ovarian cancer occur amongst women between 55 and 74 years of age. About 25% of ovarian cancer deaths occur amongst women between 35 and 54 years of age.\nThe main treatments for ovarian cancer are surgery, chemotherapy, and radiation therapy. Combinations of these treatments are used to treat ovarian cancer.\nSurgery is the usual initial treatment for women diagnosed with ovarian cancer. The ovaries, the fallopian tubes, the uterus, and the cervix are usually removed. Staging during surgery (to find out whether the cancer has spread) generally involves removing lymph nodes, samples of tissue from the diaphragm and other organs in the abdomen, and fluid from the abdomen. If the cancer has spread, the surgeon usually removes as much of the cancer as possible. This reduces the amount of cancer that will have to be treated later with chemotherapy or radiation therapy\nChemotherapy is the use of drugs to kill cancer cells. Chemotherapy may be given to destroy any cancerous cells that may remain in the body after surgery, to control tumour growth, or to relieve symptoms of the disease. Most drugs used to treat ovarian cancer are given intravenously or directly into the abdomen through a catheter.\nRadiation therapy, also called radiotherapy, involves the use of high-energy rays to kill cancer cells. Radiation therapy affects the cancer cells only in the treated area. The radiation may come from a machine or women receive a treatment called intraperitoneal therapy in which radioactive is put directly into the abdomen through a catheter.\nDeciding on a particular course of treatment is typically based on a variety of prognostic parameters and markers [Fitzgibbons, et al. (2000) Arch. Pathol. Lab. Med. 124:966-978; Hamilton and Piccart (2000) Ann. Oncol. 11:647-663)], including genetic predisposition markers BRCA-1 and BRCA-2 [Robson (2000) J. Clin. Oncol. 18:113sup-118sup].\nAlthough many ovarian cancer patients are effectively treated, the current therapies can all induce serious side effects, which diminish quality of life. Moreover, approximately 85% of the patients that have been effectively treated with platinum- and paclitaxel-based chemotherapy, including complete responses, relapse within two years after treatment.\nThe identification of novel therapeutic targets is essential for improving the current treatment of ovarian cancer patients. Recent advances in molecular medicine have increased the interest in tumour-specific cell surface antigens that could serve as targets for various immunotherapeutic or small molecule strategies.\nAmong the various elements of the immune system, T lymphocytes are probably the most adept to recognize and eliminate cells expressing foreign or tumour-associated antigens. Cytotoxic T Lymphocytes (CTLs) express the CD8 cell surface marker and are specialized at inducing lysis of the target cells with which they react via the perforin/granzyme and/or the Fas/Fas-L pathways. The T-cell receptor (TCR) for antigen of CTLs binds to a molecular complex on the surface of the target cell formed by small peptides (8-11) residues derived from processed foreign or tumour associated antigens, which associate with major histocompatibility complex (MHC) class I molecules.\nThe other major T-cell subset, helper T lymphocytes (HTLs or T helper cells), is characterized by the expression of CD4 surface marker. The T helper cells recognize slightly larger peptides (11-20 residues), also derived from foreign or tumour associated antigens, but in the context of MHC class II molecules, which are only expressed by specialized antigen presenting cells (APCs) such as B lymphocytes, macrophages and dendritic cells (DCs).\nAs a consequence of TCR stimulation of naive CTLs and HTLs by peptide/MHC complexes on APCs, the CTLs mature into effector killer cells capable of lysing (tumour) cells that express the corresponding peptide/MHC class I complex. Activated HTLs amplify CTL responses by making the APCs more effective at stimulating the naive CTLs and by producing lymphokines that stimulate the maturation and proliferation of CTLs. The potentiating effect of T helper cells occurs both in secondary lymphoid organs where the immune response is initiated and at the tumor site where CTL responses need to be sustained until the tumour cells are eliminated. Thus, one would predict that vaccines should stimulate both tumour-reactive CTLs and HTLs to generate effective antitumour immunity.\nAntigens suitable for immunotherapeutic cancer strategies should be highly expressed in cancer tissues and ideally not in normal adult tissues. Expression in tissues that are dispensable for life, however, may be acceptable.\nA number of antigens suitable for immunotherapeutic strategies in the treatment of ovarian tumors have been described so far, including MUC1, CTs, SP17 and Her2/neu.\nPolymorphic epithelial mucin (MUC1) is a transmembrane protein, present at the apical surface of glandular epithelial cells. It is often overexpressed in ovarian cancer (in more than 90% of all ovarian cancers), and typically exhibits an altered glycosylation pattern, resulting in an antigenically distinct molecule. MUC1 is in early clinical trials as a vaccine target [Gilewski, et al. (2000) Clin. Cancer Res. 6:1693-1701; Scholl, et al. (2000) J. Immunother. 23:570-580]. The tumour-expressed protein is often detectable as tumor marker in the circulation [cf. Bon, et al. (1997) Clin. Chem. 43:585-593].\nA unique class of differentiation antigens, the cancer/testis (CT) antigens, are not expressed in normal tissues except for testis and, in some cases, placenta. This fact makes CT antigens attractive targets for specific immunotherapy of cancer. The function of the majority of the CT antigens is currently unknown. Tammela et al. [Tamella, et al. (2004) Cancer Immunity 4:10-21] demonstrated that SCP-1, a CT antigen with a known role in gamete development, is expressed in 15% of ovarian cancer cases. It was suggested that because of its restricted expression in normal tissues and its aberrant expression in tumour tissues SCP-1 might serve as a potential target for vaccine therapy in ovarian cancer.\nAnother potential target for immunotherapy in patients with ovarian carcinoma is the sperm protein 17 (SP17). Sp17 was found to be expressed in the primary tumor cells from 70% of the patients with ovarian carcinoma. The restricted expression of Sp17 in normal tissue makes it an ideal target for tumour vaccine. A recombinant Sp17 protein was used with monocyte-derived dendritic cells and autologous peripheral blood mononuclear cells to generate Sp17 specific cytotoxic T-lymphocytes (CTLs). Human leukocyte antigen (HLA) class II restricted Sp17 specific CTLs were generated successfully from the peripheral blood of three patients with ovarian carcinoma at the time of disease presentation. These CTLs were able to lyse autologous Epstein-Barr virus-transformed lymphoblastoid cells in an Sp17-dependent manner. The CTLs also lysed Sp17-positive autologous tumour cells, suggesting that Sp17 is processed and presented in association with the HLA class I molecules in Sp17-positive tumour cells. [Chiriva et al. (2002) Cancer 94(9):2447-2453]\nHuman epidermal growth factor receptor 2 (Her2/neu) is an oncogene that is activated by gene amplification with the increased expression of another (normal) gene product. Her2/neu is overexpressed in 20 to 30% of patients with breast and ovarian cancer. Initial studies to develop a peptide based HER-2/neu vaccine were performed in a rat model [Disis et al. (1999) Clinical Cancer Research 5:1289-1297]. No T-cell responses or anti-body responses were observed in animals immunized with intact rat neu protein. By marked contrast, tolerance to rat neu protein in rats, could be circumvented by immunization with a peptide based vaccine. Rats immunized with neu peptides designed for eliciting CD4+ T-cell responses, generated T-cell and antibody responses specific for both the immunizing peptides and the whole protein.\nBrossart et al. demonstrated that patients with advanced breast and ovarian cancer could be efficiently vaccinated with autologues dendritic cells (DCs) pulsed with Her2/neu- or MUC1-derived peptides. In 5 out of 10 patients peptide specific CD8+ cytotoxic T lymphocytes could be detected in the peripheral blood. It was reported that MAGE-3- and CEA-peptide-specific CD8+ T cells were observed in one patient treated with MUC-1 peptide-pulsed DCs, and MUC-1 specific T-cells were observed in another patient after vaccination with HER2/neu derived peptides. It was suggested by Brossart et al. [Brossart et al. (2002) Transfus Apher Sci. 27(2):183-186] that this indicated that epitope spreading occurred in these patients upon treatment.\nEpitope spreading is a recognized phenomenon of autoimmune responses and is believed to be an exacerbating factor in CD4+ T cell-mediated autoimmune diseases. The phenomenon has been demonstrated in murine relapsing-remitting experimental autoimmune encephaleomyelitis (EAE), Theiler's murine encephalomyelitis virus-induced demyelineating disease and diabetes in the non-obese diabetic (NOD) mouse. A model has been suggested for how epitope spreading in autoimmune diseases mediated by CD4+ T cells occurs. This model is supported by direct evidence that tissue damage, TCR ligation on CD4+ T cells by MHC class II-peptide complexes, CD40-CD40 ligand interactions and CD28-mediated co-stimulation are required for epitope spreading to become manifest. It is thought that an initiating self-antigen or a persistent viral epitope, presented in MHC class II molecules on the surface of professional antigen-presenting cells (APC) residing in the target tissue, causes the activation of CD4+ T cells specific for that antigen. This T cell activation results in chronic inflammation, leading to damage of the target tissue. Tissue debris is subsequently taken up by APC which have up-regulated expression of MHC class II and co-stimulatory molecules in response to inflammatory cytokines. These APC are then capable of activating CD4+ T cells specific for secondary tissue epitopes presented by the APC. The newly activated T cells then aid in destruction of the target tissue.\nDue to the requirement for presentation by APC of exogenous antigen, epitope spreading has historically been thought of as a phenomenon unique to CD4+ T cell responses. However, recent data have indicated that cross-priming by APC can participate in the induction of CD8+ cytotoxic T lymphocyte (CTL) responses as well. In particular, bone marrow chimera studies in murine tumour models have shown that tumour-specific CTL are predominately restricted to the MHC of the host rather than that of the tumour, suggesting that indirect presentation by host APC is involved in the generation of tumour-specific CTL. Moreover, there is increasing evidence that a pathway exists whereby exogenous antigen can be presented for eventual peptide loading onto class I MHC molecules. This phenomenon is best described for dendritic cells (DC) and provides a cellular mechanism to explain the process of cross-priming. Collectively, these data suggest that it may be possible for epitope spreading to occur during a class I MHC-restricted CTL response. Because re-presentation of MHC class I-restricted tumour antigens is known to occur, it has been postulated that if tumour-bearing hosts could initiate a CTL response against a single tumour antigen, that following tumor cell damage caused by the CTL, epitope spreading might occur via a mechanism analogous to that described in CD4+ T cell-mediated autoimmune diseases. Unlike during an autoimmune response, however, CTL epitope spreading during an anti-tumour response could be beneficial to the host by possibly allowing for elimination of variant tumor cells that have lost expression of the antigen (antigen negative tumour cells). [Markiewicz et al., (2001) International Immunology 13:625-632].\nMarkiewicz et al found that immunization with the single tumour peptide P1A followed by tumour rejection led to CTL activity against a P1A− tumour, indicating that the phenomenon of epitope spreading is not limited to CD4+ T cell responses. The population of CTL included cells recognized the unrelated antigen PIE. Since this epitope was not included in the vaccine and is a mutated peptide not presented in normal tissues, the source of PIE antigen must have been the tumor cell challenge.\nBecause many patients have ovarian tumours that express neither one of the aforementioned antigens there is a need to uncover additional antigenic targets for immunotherapy to manage localized and metastatic disease. Accordingly, provided herein are molecular targets for immunotherapeutic intervention in ovarian cancers.\nThe zona pellucida (ZP) forms an extracellular glycoprotein matrix surrounding the developing and ovulated oocyte and the preimplantation embryo and is also detectable in atretic follicles. The ZP induces acrosome reaction on sperm, determines the species specificity for fertilization and prevents polyspermy in mammals. The zona pellucida contains four major glycoproteins, ZP1, ZP2, ZP3 and ZP4. In vitro studies in mice indicate that O-linked oligosaccharide side chains of ZP3 are involved in the primary binding of the sperm to the ZP3, while ZP2 contributes to the subsequent and persistent ZP binding and functions as a secondary sperm receptor.\nThe ZP glycoproteins have been studied extensively for the development of vaccines for the fertility control of animals and humans. The proposed vaccine action is the induction in female subjects of effective sustained, but reversible levels of ZP-specific antibodies that inhibit sperm-egg binding and/or prevent sperm penetration of the ZP. Passive immunization of female mice with rat monoclonal antibodies against mouse ZP2 or ZP3 resulted in localization of the antibodies to intra-ovarian oocytes and long-lasting but reversible contraception. Active immunization of female mice with ZP3-derived peptides ZP3328-342, comprising a B-cell epitope recognized by the ZP3-specific contraceptive antibody, also led to reversible albeit incomplete contraception. These ZP3 peptides also induced a T cell response to the ZP3 peptide. These CD4+ ZP3 specific T cells adoptively transfer autoimmune ovarian disease (AOD) to syngeneic recipients. Since the desired contraceptive effect of ZP3 immunization is known to be mediated by antibodies, an acceptable contraceptive ZP vaccine should induce an adequate antibody response without activation of ZP3-specific T cells. Indeed a chimeric peptide consisting of a foreign T-cell epitope from bovine ribonuclease (RNase) and a minimal and modified murine ZP3335-342 B cell epitope has been designed that elicits antibodies to ZP and has a significant contraceptive effect without causing significant oophoritis/AOD. The bovine RNase T-cell epitope stimulates helper T cell (helper T lymphocytes, HTL) responses in mice, thus potentiating the contraceptive effectiveness without inducing ZP(3)-specific T cell action and T-cell mediated ovarian damage.\nImmunisation with (self)ZP antigen has also been used to study autoimmune ovarian disease (AOD). More in particular, animal models suitable for studying AOD have been reported wherein autoimmune disease was induced using ZP antigen vaccination. For example, it was demonstrated by Rhim et al. [Rhim et al. (1992) J. Clin. Invest. 89:28-35] that in B6AF1 mice T-cell and antibody response were induced by vaccination with mouse ZP3328-342 peptide. Further studies on truncated ZP3328-342 peptides substantiated that a T cell response is sufficient for induction of oophoritis; seven of such peptides lacking antibody binding sites, elicited severe oophoritis without concomitant antibody response. These peptides include a minimal oophoritogenic peptide of eight amino acids, ZP3330-337, which overlaps the seven amino acid antibody binding site, ZP3336-342, by two residues.\nIt was reported by Bagavant et al. [Bagavant et al. (1999) Biology of Reproduction 61:635-642] that transfer of ZP3 peptide-specific T-cells into naïve recipient mice resulted in granulomatous oophoritis and enhanced ovarian expression of IL-1, TNF-α and IFN-γ. However the ovarian function of cell recipients was normal and the mice remained fertile. Antibody to ZP3 alone does not cause any ovarian pathology. Co-transfer of pathogenic T cells and ZP antibody together targets the inflammation into developing follicles leading to their destruction and the development of ovarian atrophy. In another study Bagavant et al [Bagavant et al. (2002) American Journal of Pathology 160:141-149] demonstrated that ZP3 peptide (human ZP3328-341, macaque ZP3328-341 and mouse ZP3330-342) immunization in primates can elicit a T-cell response and cause ovarian immunopathology that is similar to murine AOD.\nInternational patent application no. WO 2005/026735 (Buschmann et al.) relates to differentially expressed tumour-specific immunogenic membrane proteins and to their uses, in particular for finding at least one therapeutic molecule or compound which specifically regulates the expression of at least one of said membrane proteins, or for finding a therapeutic molecule that specifically binds to and/or interacts with any of said membrane proteins. The membrane protein can be SYPL, STOML2, RAGA, CLNS1A, PRNP, GNB2L1, GNG4, ITM2B, ITM1, TM9SF2, TM4SF6, OPRL1, LRP4, GLEPP1, TLR3 and/or ZP3. WO 2005/026735 teaches to administer the aforementioned therapeutic molecule or compound to neoplastic target cells for modulating proliferation, differentiation and/or cell migration of said neoplastic target cells. It is stated that the non-steroid dependent cancer to be treated results from the aberrant expression and/or biological activity of at least one of said immunogenic membrane proteins. It is also briefly mentioned in WO 2005/026735 that the development of a specific lesion, such as a pro-neoplastic lesion that can be found in epithelial tissues, into a neoplastic lesion can be inhibited by inoculating a subject with one of said membrane proteins adequate to produce antibody and/or T cell immune response. It is further specified that according to another embodiment the method comprises delivering one of the immunogenic membrane proteins via a vector directing expression of the said protein in vivo in order to induce such an immunological response to produce antibody to protect the subject from disease. Thus, WO 2005/026735, teaches several methods of suppressing the expression and/or biological activity of, amongst others, ZP3 membrane proteins in neoplastic target cells in order to modulate proliferation, differentiation and/or migration of said target cells, using either siRNA, receptor antagonists or antibody. WO 2005/026735 only discloses the expression of ZP3 membrane protein in certain colon cancer cells. No other reports of tumour associated expression of any ZP glycoprotein are known."}
{"text": "Traditionally, mechanical devices have been used to latch and unlatch closures such as doors, trunks, hoods, lift gates and hatches and the like in automobiles and other vehicles. It is known, however, to utilize an electro-mechanical door latch mechanism for such applications for a variety of reasons including ease of operation, lower cost and weight, improved styling opportunities, and reduced complexity. For example, a user actuated switch can be employed to trigger the release of a mechanical latch. In this regard, an electrical switch is operable to provide an input to a controller for operating the mechanical latch when the switch is actuated. In addition, modern styling and ergonomic requirements may dictate the physical configuration of the switch. For example, the switch may need to comprise an aesthetically pleasing user actuation component (e.g., a low profile button) that is of adequate size and shape so as to be easily operated by a user under a wide variety of operating conditions in a wide variety of environments.\nKnown switch technology for such applications generally incorporates a button having a first electrically conductive material comprising protrusions having the shape of “pills” or spring-like “fingers” that are insert molded or otherwise attached to the underside of the button. A second electrically conductive material comprising a set of contacts is located opposite the button on a base portion of the switch assembly. The second electrically conductive material may typically be in the form of a plate, tracks or a printed circuit board, for example. The first electrically conductive material completes a circuit in the switch when the switch is actuated by depressing the button. For example, when the button is depressed, the first conductive material bridges the contacts of the second electrically conductive material thereby closing an electric circuit.\nIn one such known switch assembly configuration, the switch assembly is also sealed from the atmosphere. During its manufacture, a fixed volume air is captured in the space between the button and the base portion of the switch assembly. As such, when the ambient temperature of the switch assembly changes, so too does the volume of the air trapped within the switch assembly. Under hotter ambient temperature conditions the volume of air within the switch assembly expands; under colder ambient temperature conditions, the volume of air within the switch assembly contracts.\nIn such a design, changes in the switch assembly's operating environment, such as extreme changes in ambient temperature, for example, can impact the perceived operation of the switch to a user. For example, the switch assembly may not reliably provide satisfactory and perceptible tactile feedback to the user signifying actuation of the switch. In such a case, depression of the button may instead provide an unsatisfactory continuous resistance to the user causing the user to be unsure whether the switch has been properly actuated.\nConsequently, it is desirable to provide a switch assembly having a reliable and cost-effective actuation mechanism that also provides satisfactory tactile feedback to a user for signifying proper actuation of the switch."}
{"text": "1. Field of the Invention\nThe subject invention applies to the provision and display of video information and has particular application to apparatus for displaying recorded video programs, although not being limited thereto.\n2. Description of the Prior Art\nThe detrimental effect of flutter on the display of video signals is well known, as is the fact that video information recording and playback equipment is a typical, although by no means exclusive, source of flutter.\nThe thrust of prior-art endeavor has been to improve the equipment responsible for the flutter. For instance, design and precision of video recording and playback system, including video tape recorders, have been continuously improved toward a reduction of flutter to tolerable proportions. This has neither been consistently successful nor has it helped in realistically bringing the equipment in question within the reach of the public or of educational and other institutions on a large scale. Rather, the requisite improvements for an implementation of prior-art proposals in this area have greatly increased the complexity and expense of manufacture and design of the apparatus in question.\nAn aggrevating factor in this area has been the tendency of conventional television display sets to amplify the effect of flutter in received signals. The source of this problem has been traced to a phase lock loop or electronic flywheel circuit present in most sets for stabilizing synchronizing signals against noise above a certain frequency. In general, prior-art stabilization circuits of this type have the advantage of permitting the synchronization drive to adjust to synchronization frequency differences when switching stations or to drift occurring in any station, while at the same time rendering the synchronization drive relatively immune to frequently occurring noise disturbances.\nHowever, it will now be shown that these otherwise beneficial stabilization circuits amplify the effect of flutter in certain instances. Ideally, the horizontal position of a picture element in a video display is related to the phasing of the corresponding horizontal sync pulse. In practice, this is no longer the case if the horizontal sync has been stabilized by the stabilizing circuit of the set to a relatively steady time averaged value.\nFor instance, the line rate as delivered by a video tape recorder in the playback mode will have a flutter component. If .nu..sub.o is the nominal horizontal sync frequency and A is the peak amplitude of a flutter component of frequency .apprch., then the delivered line rate .nu. will be EQU .nu.=.nu..sub.o [1+A sin 2.pi..psi.t] (1)\nThe horizontal position of a picture element ideally will be at some fixed phase with respect to the beginning of the particular line. If .phi..sub.h is the absolute phasing of the line beginning, we have ##EQU1## so that .phi..sub.h is EQU .phi..sub.h =2.pi..intg..nu..sub.o [1+A sin 2.pi..psi.t]dt (3)\nSince we are interested in the excursion of the variational component due to flutter and the peak amplitude of the excursion, the integral is taken between EQU t =0 and t =1/4.psi. (4)\nvariational component is ##EQU2## or EQU .phi..sub.h =.nu..sub.o A/.psi. (6)\nas the amplitude excursion of the screen element when the sync is stabilized in the television display set.\nSince the horizontal sweep is linear with time, this excursion expressed as a fraction of screen width is ##EQU3##\nwherein .GAMMA. is the peak amplitude of the excursion of a picture element and A is the peak amplitude of the flutter component whose frequency is .psi..\nAs to the effect of flutter on the image display, we are concerned with the ratio ##EQU4##\nFor a nominal horizontal sync frequency of .nu..sub.o =15, 750 Hz and a flutter component frequency of .psi.=100 Hz, we obtain a ratio .GAMMA./A of ##EQU5##\nVideo recording and playback machines in which all flutter components remain below 0.1% or A A.ltoreq. 0.001 are expensive and complex. But even at that low value, the effects of flutter due to the above mentioned stabilizing action is still very objectionable.\nFor instance, if we take the above example of equation (9), we find that the picture element would execute excursions with an amplitude of EQU 25.sup.. 10.sup..sup.-3 =0.025 (10)\nas a fraction of screen width.\nFor a screen width of about 20 inches (508 millimeters), for instance, the excursion amplitude of picture elements would be about 0.5 inches (12.7 millimeters) or 1 inch (25.4 millimeters) peak-to-peak.\nSuch a deviation is intolerable on its face. To reduce such excursions to, say, 0.04 inches or about one millimeter under the permises of our above example, it would be necessary to reduce the flutter of the video recording and playback machine to about 4.10.sup..sup.-5 or 0.004%!\nThe general equation applying to these cases is ##EQU6## wherein: A is the peak amplitude of flutter;\n.OMEGA. is the peak-to-peak excursion of the picture element on the display screen; PA1 .nu..sub.o is the nominal horizontal sync frequency; and PA1 .lambda. is the screen width.\nThe resulting precision requirements are far too stringent in most practical applications.\nIn their quest for a workable solution, proposals have been made which necessitate an invasion and modification of the existing circuitry of the television display set. With respect to such proposals, most experts in the field agree that systems which are inherently incapable of operation without invasion and modification of the existing receiver circuitry have practically no chance of widespread acceptance among television receiver owners."}
{"text": "The (+)-threo acid is an important raw material of the therapeutically active (2S,3S)-3-acetoxy-5-dimethylaminoethyl-2-(4-methoxyphenyl)-2,3-dihydro-1,5 -benzothiazepin-4(5H)-one hydrochloride (generic name: diltiazem), a known antianginal drug acting as a calcium antagonist.\nThere are several processes known from the literature for the preparation of the (+)-threo acid.\nA resolution process using L-lysine as resolving agent is described in the published German patent application (DE-OS) No. 3,337,176. However, the particularly costly L-lysine is used in a 4-fold excess and,; in addition, it is difficult and expensive to recover L-lysine in its acid form from the mother liquors.\nAccording to another method [Helv. Chim. Acta 67, 916 (1984)], the racemic threo acid is resolved by using cinchonidine in ethanol. This method is characterized by the use of an extraordinarily high amount of alcohol as solvent and by a 48-hour time demand of crystallization.\nAccording to the method described in the U.S. Pat. No. 4,416,819, an equivalent amount of (+)-.alpha.-phenylethylamine each is employed for both forms of the racemic acid in order to resolve the racemic threo acid in water. In our own examinations, a (+)-threo acid with an optical purity of only 87%, [.alpha.].sub.D.sup.20 =+302.degree. (c=0.3, ethanol), was obtained in a yield of 79% when the resolution and subsequent crystallization of the diastereomeric salt were carried out by using the method cited above.\nIn the process described in the Hungarian patent specification No. 193,230, the acid addition salt of 0.5 to 0.6 molar equivalent of (+)-.alpha.-phenylethylamine is used as resolving agent; a method, useful to enrich the valuable (+)-threo acid in the remaining threo acid mixture is also described in the above specification."}
{"text": "This application is based on application No. H11-008674 filed in Japan on Jan. 18, 1999, the entire content of which is hereby incorporated by reference.\n1. Field of the Invention\nThe present invention relates to an optical system for a projector, and particularly to such an optical system for a projector as incorporates a reflection-type spatial light modulator, such as a digital mirror device(trademark) (or DMD(trademark) for short, manufactured by Texas Instruments Incorporated; hereinafter referred to simply as a digital mirror device or DMD), that is provided with a large number of variable-reflection-angle mirror elements that can vary the reflection angle of the light incident thereon in accordance with a video signal so that only the light corresponding to the video signal will be reflected toward a projection optical system.\n2. Description of the Prior Art\nIn recent years, as higher and higher resolution is desired in images in general, also in the field of projectors, development of techniques has been sought that achieve a substantial increase in the number of pixels without making the optical system larger. One attempt to meet such needs is the development of a projector employing a DMD.\nA DMD is produced by forming a large number of minute rectangular high-reflectance mirror elements, of which the inclination is variable in accordance with a video signal, on a silicon memory chip by a process as used to manufacture a semiconductor device. A projector employing such a DMD, by varying the inclination of those mirror elements, controls the direction in which it reflects the illumination light incident thereon in such a way as to converge only the desired parts of the reflected light on a screen and thereby project a desired image thereon.\nIn accordance with a video signal, the individual mirror elements of a DMD, when in an on state, reflect light toward a projection optical system and, when in an off state, change their inclination to reflect light in a different direction so that the light will not enter the projection optical system. However, owing to restrictions imposed on the design of a projector, the light reflected from the mirror elements when they are in an off state is shone on a side-wall portion of a prism constituting a prism system disposed between the DMD and the projection optical system, and therefore there is a possibility of this light being further reflected therefrom so as to enter the projection optical system.\nIn particular, when the prism is surrounded by a medium, such as air, having a small refractive index, most of the above-mentioned light shone on the side-wall portion of the prism is reflected from the inner surface of the side-wall portion of the prism, and this greatly increases the amount of such secondary-reflection light. This secondary-reflection light (unnecessary, or stray, light), when it enters the projection optical system, may cause another image (a ghost) separate from the normal image to appear on the screen.\nOne way to prevent such entry of secondary-reflection light into the projection optical system is to form the side-wall portion of the prism into a diffusive surface that diffuses the light incident thereon. Another way is to paint the outer surface of the side-wall portion of the prism black, or to vapor-deposit on that surface a light-absorbing dielectric film, or to affix to that surface a member having a bottom surface so shaped as to absorb light so that this bottom surface will absorb light and convert it into heat.\nThe method of absorbing light and converting it into heat by the use of a member having a bottom surface so shaped as to absorb light is adopted, for example, in the projector-oriented optical system disclosed in Japanese Laid-Open Patent Application H9-96867. In this optical system, a heat-dissipating member having a comb-tooth-shaped bottom surface is arranged so as to face a side wall of a prism, with a shock-absorbing pad in between, in such a way that this bottom surface is kept in close contact with the shock-absorbing pad. In this optical system, the recessed portions of the comb-tooth-shaped bottom surface absorb light and convert it into heat, and the heat-dissipating member as a whole dissipates the resulting heat.\nHowever, even if light is diffused, it is inevitable that part of the diffused light will enter the projection optical system. It is practically impossible to form a black thin film that completely absorbs the light incident thereon, and therefore, even if a black thin film is applied, some light, left unabsorbed, may enter the projection optical system. It is difficult to produce a member having so intricate a shape as to absorb completely the light incident thereon, and therefore it is inevitable that part of the incident light will be reflected or diffused in the vicinity of the protruding portions of the comb-tooth-shaped bottom surface of this member.\nAn object of the present invention is to provide a projector-oriented optical system in which entry of unnecessary light into a projection optical system is prevented more strictly than ever.\nTo achieve the above object, according to one aspect of the present invention, an optical apparatus is provided with: a reflection-type spatial light modulator having a plurality of minute variable-reflection-angle mirrors that individually deflect the light incident thereon in one of two different directions, namely in a first direction or in a second direction, in accordance with a signal fed in; an optical system to which the light deflected in the first direction by the modulator is directed; a prism disposed between the modulator and the optical system so as to direct the light deflected in the first direction by the modulator to the optical system and direct the light deflected in the second direction by the modulator to a side face of the prism; and a light-absorbing member shaped like a plane-parallel plate and disposed with an entrance face thereof kept in close contact with the side face of the prism. The light-absorbing member is so designed that the light deflected in the second direction by the modulator enters the light-absorbing member through the entrance face thereof but does not exit from the light-absorbing member through the entrance face back into the prism.\nIn this optical apparatus, preferably, the light-absorbing member is so designed that the distance t between the entrance face and a second flat face opposite to the entrance face fulfills the following condition:\ntxe2x89xa7[1xe2x88x92{(n1/n2)sin xcex11}2]xc2xd/2K\nwhere n1 represents the refractive index of the prism, n2 represents the refractive index of the light-absorbing member, xcex11 represents the minimum angle of incidence of light incident on the entrance face, and K represents the absorption coefficient of the light-absorbing member.\nAccording to another aspect of the present invention, a projector is provided with: an illumination optical system that includes a light source and that emits the illumination light generated by the light source; a reflection-type spatial light modulator having a plurality of minute variable-reflection-angle mirrors that individually deflect the illumination light emitted from the illumination optical system in one of two different directions, namely in a first direction or in a second direction, in accordance with a signal fed in; a projection optical system to which the light deflected in the first direction by the modulator is directed; a prism disposed between the modulator and the optical system so as to direct the illumination light emitted from the illumination optical system to the modulator, direct the light deflected in the first direction by the modulator to the optical system, and direct the light deflected in the second direction by the modulator to a side face of the prism; and a light-absorbing member shaped like a plane-parallel plate and disposed with an entrance face thereof kept in close contact with the side face of the prism. The light-absorbing member is so designed that the light deflected in the second direction by the modulator enters the light-absorbing member through the entrance face thereof but does not exit from the light-absorbing member through the entrance face back into the prism.\nIn this projector, preferably, the light-absorbing member is so designed that the distance t between the entrance face and a second flat face opposite to the entrance face fulfills the following condition:\ntxe2x89xa7[1xe2x88x92{(n1/n2)sin xcex11}2]xc2xd/2K\nwhere n1 represents the refractive index of the prism, n2 represents the refractive index of the light-absorbing member, xcex11 represents the minimum angle of incidence of light incident on the entrance face, and K represents the absorption coefficient of the light-absorbing member.\nIn the optical apparatus and the projector described above, the light shone on the side face of the prism enters the light-absorbing member. The light having entered the member is either absorbed by the member and thereby converted into heat or, if absorbed incompletely, driven out of the member. Here, the light driven out of the member is never driven out of the member through the face that is kept in close contact with the prism. Specifically, for example, the light-absorbing member is so designed that the light having entered the member is completely absorbed within the member before striking again the face that is kept in close contact with the prism. In this way, it is possible to prevent unnecessary light from returning to the prism (and thus entering the (projection) optical system).\nAssume that FIG. 1 is a sectional view of the light-absorbing member used in the optical apparatus or projector described above along the plane that is perpendicular to both the first and second faces (what is called respectively the entrance face and the second flat face above) of the member. In FIG. 1, consider, for example, a case where light, coming from inside the prism 22, enters the light-absorbing member 5 through the bottom face 5a thereof (i.e. the first face that is kept in close contact with the prism 22), at a position P1 thereon, is then reflected from the top face (the second face) 5b, at a position P2 thereon, and then travels toward the bottom face 5a, toward a position P3 thereon. Here, assume that the angle of incidence and the angle of emergence at the position P1 are xcex11 and xcex12, respectively, the distance between the bottom and top faces 5a and 5b (i.e. the thickness of the light-absorbing member 5) is t, the absorption coefficient of the light-absorbing member 5 is K, and the distance between P1 and P2, which equals the distance between P2 and P3, is L.\nIf the light that has entered the light-absorbing member 5 at the position P1 is absorbed before reaching the point P3 (i.e. while traveling an optical path length 2L), then\n1xe2x89xa62LKxe2x80x83xe2x80x83(1)\nHere, according to the trigonometric function,\nxe2x80x83L=t/cos xcex12xe2x80x83xe2x80x83(2)\nAccording to Snell\"\"s law,\nn1 sin xcex11=n2 sin xcex12\nsin xcex12=(n1/n2)sin xcex11xe2x80x83xe2x80x83(3)\nCombining Expressions (2) and (3) gives\nL=t/[1xe2x88x92{(n1/n2)sin xcex11}2]xc2xdxe2x80x83xe2x80x83(4)\nSubstituting Expression (4) in Expression (1) gives\n1xe2x89xa62Kt/[1xe2x88x92{(n1/n2)sin xcex11}2]xc2xd\ntxe2x89xa7[1xe2x88x92{(n1/n2)sin xcex11}2]xc2xd/2Kxe2x80x83xe2x80x83(5)\nAs described above, for the light that has entered the light-absorbing member 5 at the position P1 to be absorbed before reaching the point P3 (i.e. while traveling an optical path length 2L), Expression (5) needs to be fulfilled. Expression (5) defines the lower limit of the thickness t of the light-absorbing member 5. Expression (5) shows that, the smaller the angle of incidence xcex11, the greater the lower limit of the thickness t. Accordingly, as long as the light-absorbing member 5 is so formed as to have a thickness t that fulfills Expression (5), assuming that the minimum angle of incidence of the light entering the member 5 through the bottom face 5a thereof equals xcex11, light never returns to the prism through the bottom face 5a."}
{"text": "This invention relates to a custom-formable knee immobilizer product, knee immobilizer and method.\nThe invention has particular application in the orthopedic medical field, where injured or diseased knees are often supported and immobilized with a splint or brace so that the patient can continue to walk, albeit with a stiff leg, while the injury heals. Typical injuries which can sometimes require immobilization may include, but are not limited to, pulled or torn quadriceps muscles and tendons and/or the patellar tendon, chondromalacia patellae, pulled or torn medial or lateral collateral ligaments, and posterior or anterior cruciate ligaments. Certain related degenerative conditions, as well post-operative recovery and therapy following a wide range of orthopedic surgical procedures, may also require immobilization of the knee.\nIn such cases, it is essential to immobilize the knee in such a way as to prevent articulation., which can strain healing muscles, tendons, ligaments and bone surfaces. In addition, the knee must be stabilized to prevent eversion and inversion.\nThe invention of the present application takes advantage of polymer chemistry to permit quick and easy molding of an immobilizer to the knee and adjacent thigh and calf portions of the leg. The close, custom fit achieved is in distinct contrast to so-called xe2x80x9cone size fits allxe2x80x9d braces or splints wherein bendable plastic rods or stays provide support, and are held in place by a heavy fabric tube.\nThe particular embodiment of the immobilizer disclosed in this application will accommodate a wide range of sizes and can be fitted to either the right or left knee. Thus, a much reduced inventory is required. This feature also substantially reduces design and manufacturing costs, and promotes use through ease of fitting. The custom-fit of the immobilizer not only permits easy removal for bathing, dressing or adjustment, but also easy and mistake-proof replacement even by children.\nPrior art immobilizers include numerous types of splints and braces which typically include a soft component to place near the skin and a semi-rigid metal or plastic stays which are bent into a shape conforming to the aspect of the knee against which the stay resides. The soft component, for example, fiber padding, foam or an air bladder, is intended not only to provide a cushion, but also to accommodate itself to the varying configurations of differing sized and shaped body parts. For this reason, the cushioned part is substantially greater in thickness than required merely to provide the protection from injury or irritation caused by direct contact with the stays.\nSome prior art knee braces include stays which are constructed of thermosetting materials which are heated and then formed to the body while heated. These products require a source of heat, and are susceptible to either over-or-underheating. In addition, body heat itself can soften or increase the flexibility of the stays, thereby decreasing the effectiveness of the protection offered.\nThe present invention permits quick and easy application of a immobilizer to a knee in such a way as to achieve a true custom fit. The. moisture curable resin system used results in a very rigid splint which holds the shape of the molded splint to a very high degree. No heat is required, and a source of water is the only additional material necessary to achieve a cure. Atmospheric moisture alone will cure the splint into its hardened position in a relatively short period of time, but in practice the resin in or on the splint will typically be activated by dipping the splint in water and then removing the excess water by rolling the splint in a towel immediately before the splint is applied to the knee.\nTherefore, it is an object of the invention to provide a custom-formable knee immobilizer.\nIt is another object of the invention to provide a knee immobilizer which can be molded to a knee to stabilize the knee against articulation, eversion and inversion.\nIt is another object of the invention to provide a knee immobilizer which can be custom-formed to a particular patient.\nIt is another object of the invention to provide knee immobilizer having a shape which permits it to be used on either the left or the right knee.\nIt is another object of the invention to provide knee immobilizer which hardens in the presence of moisture to form a very rigid but very lightweight splint portion.\nIt is another object of the invention to provide a knee immobilizer which is stored in a moisture-proof pouch until ready for being custom-fitted to the knee.\nIt is another object of the invention to provide a knee immobilizer which is custom-molded to a patient\"\"s knee in such a way that it can be initially held in place with an elongate elastic strip to reduce edema, and thereafter held in place with a strap or straps during later stages of healing.\nThese and other objects of the present invention are achieved in the preferred embodiments disclosed below by providing a custom-formable knee immobilizer product for being custom-formed to the shape of a knee while flexible and upon hardening providing a rigid, knee-supporting custom fit. The knee immobilizer product comprises an outer container formed of moisture-impervious material and a plurality of elongate knee splint segments positioned in the container in substantially moisture-free conditions and sealed therein against entry of moisture until use. Each of the knee splints comprises an elongate substrate, and a reactive system impregnated into or coated onto the substrate. The system remains stable when maintained in substantially moisture-free conditions and hardens upon exposure to moisture to form a rigid, self supporting structure. A flexible protective outer cover is positioned over the substrate along its length to enclose the substrate. The substrate and outer cover thus form a unitary structure for being molded while flexible to an aspect of the knee and adjacent thigh and calf portions of the leg. A bandage is provided for being positioned on and held in a closely-conforming configuration around the knee and adjacent thigh and calf portions of the leg. Splint supports cooperate with the splints and the bandage for maintaining the plurality of splints on the bandage in closely conforming position in vertically spaced-apart relation around the knee for substantially immobilizing the knee in a predetermined recovery position.\nAccording to one preferred embodiment oft he invention, the splint supports include first and second complementary hook-and-loop elements on respective mating surfaces of the bandage and the splints.\nAccording to another preferred embodiment of the invention, the splint supports include a plurality of vertically-extending pockets formed in the bandage for positioning and supporting the vertically-extending splints around the knee.\nAccording to yet another preferred embodiment of the invention, the container is fabricated of an aluminum foil laminate having an outer tear resistant layer, a central aluminum foil layer and an inner heat sealable plastic layer.\nAccording to yet another preferred embodiment oft he invention, the substrate is formed from a plurality of knitted or woven fabric layers.\nAccording to yet another preferred embodiment of the invention, the bandage is formed from a foam sheet material.\nPreferably, the foam material is chosen from the group consisting of open or closed cell EVA or polyurethane.\nAccording to yet another preferred embodiment of the invention, the bandage includes a fabric material covering at least one side of the foam sheet material.\nAccording to yet another preferred embodiment of the invention, the fabric of the outer cover is chosen from the group consisting of woven, knitted or non-woven fabric.\nAccording to yet another preferred embodiment of the invention, the reactive system comprises a blended polyisocyanate, polyol, catalyst and stabilizer.\nAccording to yet another preferred embodiment of the invention, the bandage is a tube for being pulled onto the leg over the foot.\nAccording to yet another preferred embodiment of the invention, the bandage includes a flat sheet having cooperating mechanical fasteners for releasably holding the sheet around the knee.\nAccording to yet another preferred embodiment of the invention, the mechanical fasteners comprise first and second complementary hook-and-loop elements on mating surfaces of the flat sheet.\nAccording to yet another preferred embodiment of the invention, an elongate elastic strip is provided for being wrapped around the knee over the bandage and splints during curing of the moisture-curable resin and thereby holding the bandage and splints in closely conforming position on the knee.\nAccording to yet another preferred embodiment of the invention, the splints are approximately one inch wide.\nAccording to yet another preferred embodiment of the invention, at least one retaining strap is provided for releasably encircling and maintaining the custom-molded bandage and splints in supporting position on the knee.\nIn accordance with a subcombination of the product, a custom-formable knee immobilizer is provided for being custom-formed to the shape of a knee while flexible and upon hardening providing a rigid, knee-supporting custom fit. The knee immobilizer comprises a plurality of elongate knee splints, each comprising an elongate substrate and a reactive system impregnated into or coated onto the substrate. The system remains stable when maintained in substantially moisture-free conditions and hardens upon exposure to moisture to form a rigid, self supporting structure. A flexible protective outer cover is positioned over the substrate along its length to enclose the substrate. The substrate and outer cover form a unitary structure for being molded while flexible to an aspect of the knee and adjacent thigh and calf portions of the leg. A bandage is positioned on and held in a closely-conforming configuration around the knee and adjacent thigh and calf portions of the leg. Splint supports cooperate with the splints and the bandage for maintaining the plurality of splints on the bandage in closely conforming position in vertically spaced-apart relation around the knee for substantially immobilizing the knee in a predetermined recovery position.\nAn embodiment of the method of therapeutically immobilizing a knee according to the invention includes the steps of providing a plurality of elongate knee splints, each comprising an elongate substrate and a reactive system impregnated into or coated onto the substrate. The system remains stable when maintained in substantially moisture-free conditions and hardening upon exposure to moisture to form a rigid, self supporting structure. A flexible protective outer cover is positioned over the substrate along its length to enclose the substrate, the substrate and outer cover forming a unitary structure for being molded while flexible to an aspect of the knee and adjacent thigh and calf portions of the leg. A bandage is provided for being positioned on and held in a closely-conforming configuration around the knee and adjacent thigh and calf portions of the leg. Splint supports cooperate with the splints and the bandage for maintaining the plurality of splints on the bandage in closely conforming position in vertically spaced-apart relation around the knee for substantially immobilizing the knee in a predetermined recovery position. The bandage is placed in a closely-conforming configuration around the knee and adjacent thigh and calf portions of the leg. The splints are wetted and placed into engagement with the bandage by cooperation with the splint supports in vertically spaced-apart relation around the knee to thereby substantially immobilize the knee in a predetermined recovery position. The bandage and engaged splints are held in the closely-conforming configuration for a period of time sufficient to allow the resin to harden the splints.\nAccording to one preferred embodiment of the invention, the step of placing the bandage in a closely-conforming configuration around the knee and adjacent thigh and calf portions of the leg comprises the step of pulling a tubular bandage over the foot and onto the knee portion of the leg.\nAccording to another preferred embodiment of the invention, the step of placing the bandage in a closely-conforming configuration around the knee and adjacent thigh and calf portions of the leg comprises the step of placing a flat sheet around the knee and releasably securing the sheet around the knee with mechanical fasteners.\nAccording to yet another preferred embodiment of the invention, the step of securing the sheet around the knee with mechanical fasteners comprises the step of connecting first and second complementary hook-and-loop elements carried on matable surfaces of the flat sheet.\nAccording to yet another preferred embodiment of the invention, the method includes the step of maintaining the bandage and engaged splints in the closely-conforming configuration for a period of time sufficient to allow the resin to harden the splints comprises the step of wrapping an elongate elastic strip around the knee. over the bandage and splints during curing of the moisture-curable resin, thereby holding the bandage and splints in closely conforming position on the knee.\nAccording to yet another preferred embodiment of the invention, the step of placing the splints into engagement with the bandage comprises the step of inserting the splints into respective vertically-oriented pockets formed in the bandage.\nAccording to yet another preferred embodiment of the invention, the step of placing the splints into engagement with the bandage comprises the step of attaching to the bandage with complementary hook-and-loop elements."}
{"text": "This invention relates generally to medical devices and more particularly to hemostatic devices.\nAs will be appreciated by those skilled in the art various surgical procedures are now being carried out intravascularly or intralumenally. For example in the treatment of vascular disease, such as atherosclerosis, it is a common practice to invade the artery to insert an instrument, e.g., a balloon or other type of catheter to carry out the procedure within the artery. Such procedures usually involve the percutaneous puncture of the artery so that an introducer sheath can be inserted into the artery and thereafter the instrument, e.g., catheter, itself can be inserted through the sheath to the operative position within the artery. Such procedures unavoidably present the problem of stopping the bleeding at the percutaneous puncture after the procedure has been completed and after the instrument (and any introducer sheaths used therewith) have been removed. At present such bleeding is stopped by the application of direct digital pressure over the puncture site by a trained physician or other suitably trained medical personnel. Such direct pressure has to be applied for a sufficiently long time for hemostasis to occur so that the opening is effectively closed against further bleeding. In the case of punctures into femoral or superficial femoral arteries the pressure may have to be applied for as long as forty-five minutes for hemostasis to occur. Not only is this direct digital pressure application procedure wasteful of time by highly skilled medical professionals, the procedure results in a substantial reduction, if not virtual arrest, of the flow of blood through the vessel. Since thrombosis is one of the major calamities that can occur in the immediate post operative period, any reduction in blood flow, such as caused by the application of digital pressure, is undesirable.\nApplicator devices have been disclosed in the patent literature for inserting an absorbent plug or member into the vagina. Such devices basically comprises a tubular element adapted to be inserted into the vagina and having a plug of absorbent material located therein. The device also includes a plunger to push the plug out of the tubular element into the vagina. The plug also includes a thread or string attached to it to enable the plug to be retrieved from the vagina. Examples of such devices are shown in U.S. Pat. Nos. 1,191,736 (Roberson) and 1,794,221 (Washburn et al.).\nWhile such devices are suitable for their intended purposes, there is no suggestion of their use, nor are they suitable for insertion into an opening in the wall of a blood vessel or other bodily lumen or duct to seal that opening.\nThe patent literature also includes devices for closing an opening in a blood vessel using sutures, see U.S. Pat. No. 4,587,969 (Gillis). Other means and techniques for closing a wound are disclosed in U.S. Pat. No. 4,606,337 (Zimmermann et al.).\nNone of the prior art teaches the use of simple means for effecting the closure of an opening, e.g., puncture, in the wall of a blood vessel, duct or lumen, by plugging the opening and without requiring sutures or the application of digital pressure."}
{"text": "1. Field of The Invention\nThis invention relates to a vacuum vessel having a member cooled to a low temperature in use, for example, to the temperature of liquid helium, and means for cooling the cooled member. More particularly, but not exclusively, the present invention relates to a cryopump for condensing and adsorbing gas molecules on a cryopanel surface cooled to an ultra-low temperature, such as is used in a nuclear fusion reactor, etc.\n2. Description of The Prior Art\nThe article \"Cryopumps and cryogenic systems of prototype injector unit for JT-60\", by T. Shibata et al., 9th International Cryogenic Engineering Conference, May 1982, describes a large scale cryogenic apparatus. A cryopump unit is used in each of a plurality of neutron beam injection devices disposed around a nuclear fusion apparatus. Each cryopump unit is cooled by liquid helium and liquid nitrogen. The cryopanel of each cryopump which is used for condensing and thus adsorbing gas molecules and exhausting them at a high speed is cooled to about 3.7K by liquid helium. A heat shield plate cooled to a low temperature of about 80K by liquid nitrogen is disposed around the cryopanel of each cryopump unit to avoid the cryopanel being directly heated by radiant heat from normal and high temperature portions outside the pump.\nIn known arrangements, the liquid helium that cools the cryopanel of each cryopump unit is produced collectively by a common large-sized liquid helium liquefaction refrigerator and is transferred to each cryopump unit through a heat insulated piping. Low temperature helium gas evaporating inside the cryopanel is recovered by the same large-scale liquid helium liquefaction refrigerator through another heat insulated piping. Liquid nitrogen, too, is produced collectively by a common large-sized liquid nitrogen liquefaction machine or is transferred from the same liquid nitrogen storage tank to each cryopump unit through a heat insulated piping, and the low temperature nitrogen gas evaporating inside the heat shield plate is emitted to the atmosphere.\nIn this known, large-sized cryopump, liquid helium is transferred from the large-scale liquid helium liquefaction refrigerator to each cryopump unit through a long heat insulated piping, and a large quantity of liquid helium evaporates inside the piping. This evaporated low temperature helium gas cannot be utilized for cooling the cryopanel, and moreover, increases the fluidization pressure loss inside the piping. Therefore, the pipe diameter of the heat insulating piping must be increased.\nHeat entering the insulated piping is generally about 1 W/m, and when the length of the heat insulated piping reaches 200 m, liquid helium corresponding to heat of 200 W evaporates. Heat entering the helium at the cryopump is generally from several to some dozens of watts, so that most of the liquid helium is consumed merely for cooling the heat insulated piping.\nThe known large-sized cryopump thus needs a large quantity of liquid helium for cooling the heat insulated piping. For this reason, a large-sized liquid helium liquefaction apparatus having a large power consumption has been required.\nIf gas leaks into the vacuum region of the vacuum insulated piping, the heat insulating effect is largely destroyed and a large quantity of liquid helium inside the piping evaporates, so that the heat transfer function stops, cooling the piping from ambient temperature cannot be achieved, and the function of the cryopump often stops.\nJapanese Utility Model publication no. JP-U-62-167875 discloses cryopanels cooled by a helium circuit including a liquefying refrigerator and a liquid helium tank. A heat shield or baffle of louvers protects a cryopanel from ambient temperature parts.\nThe prior art technique described above cools a group of cryopumps with liquid helium produced by a concentrated type large-sized liquid helium liquefaction machine and transferred through heat insulated pipings disposed in the ambient atmospheric space. It therefore involves the following problems:\n(1) high electric power is necessary for cooling the cryopanels, PA0 (2) an expensive large-sized liquid helium liquefaction apparatus is necessary, PA0 (3) long heat insulated piping having large diameters are necessary, PA0 (4) since the long piping must be first cooled, a long time (several days) is necessary for cooling the cryopanel, PA0 (5) heating the cryopanel needs a long time (several days), likewise. PA0 (6) the reliability of the cooling system is low."}
{"text": "To manage today's tremendous data handling requirements, most enterprises employ some form of a data facility. Typical data facilities house data storage and processing equipment that can perform a variety of data storage and computational tasks. Data facilities may also host web servers, Internet services, and other enterprise-based services. While maintaining a data facility enables an enterprise to manage large quantities of data, the energy costs associated with running a data facility present a significant expenditure. A large percentage of data facility energy costs arise from the environmental controls required to ensure that the environment within the data facility is maintained within suitable parameters for the equipment contained in the facility. Examples of the environmental controls include cooling, air flow, humidity controls, power regulators, and so on. All of these controls work together to attempt to create an environment in which the data facility equipment can operate at maximum efficiency and thus decrease the overall energy costs of the data facility.\nProblems can arise in a data facility when environmental conditions fail to remain within acceptable operating parameters. Hot spots can cause equipment to run at less than optimal efficiency and at extremes can result in equipment failure and service interruptions. Excess humidity can allow condensation to form in and around data facility equipment and result in data processing and storage errors and ultimately, equipment failure. To control environmental conditions such as temperature and humidity, a data facility administrator needs to be aware of both global and local environmental conditions within the facility.\nTo enable data facility designers and administrators to determine optimal placement and settings for environmental controls, some form of environmental monitoring is desirable. Most current forms of environmental monitoring are difficult to implement and tend to create an incomplete and inaccurate image of data facility environmental conditions."}
{"text": "This invention relates to small electric motors of the type used in axial flow air impeller units commonly found in computers, copy machines etc. More particularly, the invention relates to an improved bearing and seal assembly for use about a central shaft and within a stator or rotor of such a motor. The bearing and seal assembly in a motor of this type must be of a relatively simple and inexpensive construction and yet highly efficient and dependable in use throughout the substantial life span of the unit.\nIn U.S. Pat. No. 3,614,261 entitled BEARING SEAL MEANS FOR AN ELECTRICALLY DRIVEN FAN issued Oct. 19, 1971 to Karl E. Friese et al, a bearing and seal assembly is disclosed with a sleeve bearing disposed within a ferrule or sleeve which provides a capillary sealing function at one end and which is enlarged and closed at an opposite end to provide a lubricant chamber about the bearing. The ferrule or sleeve is, however, of relatively heavy and expensive construction and the capillary sealing function is provided by a cup-shaped projecting end portion of the sleeve and an associated annular felt ring. Close tolerances and precise concentricity or coaxial relationship between the sleeve bearing and the portion of the sleeve providing the capillary seal are of course desired but have been found somewhat difficult to achieve and the construction has also been found deficient in this regard.\nIt is the general object of the present invention to provide a bearing and seal assembly which includes a bearing and seal confining sleeve and which provides for the necessary close tolerances, precise concentricity or coaxial relationship of elements, and which is yet of a desirably simple and inexpensive construction yielding a high degree of efficiency and dependability in use.\nA further object of the invention resides in the provision of an assembly method whereby the sleeve receives the bearing and seal in press fit engagement and precisely centers the elements about an assembly pin with the sleeve press fit or adhesively secured within the bore of a motor stator or rotor."}
{"text": "1. Field of the Invention\nThe present invention relates to an organic electroluminescence (EL) panel having organic EL elements arranged in a matrix on a substrate, and more particularly to such an organic EL panel which can effectively prevent moisture from entering the organic EL elements.\n2. Description of Related Art\nOrganic EL display panels have conventionally been known as one example of flat display panels. Because, unlike liquid crystal display (LCD) panels, organic EL display panels are self emissive and because organic EL display panels are bright and clear flat display panels, their widespread use is highly expected.\nAn organic EL display comprises a great number of organic EL elements arranged in a matrix, and employs these organic EL elements as pixels of a display. Such organic EL displays can be driven passively or actively, similar to LCDs, and, also as is the case with LCDS, active matrix displays are more preferable. More specifically, in active matrix displays, switching elements (typically two elements including a switching element and a driving element) are provided for each pixel and display of each pixel is controlled by controlling the switching elements, whereas in passive driving, a switching element is not provided for each pixel. Of these two types of displays, the active matrix is more preferable because much more precise displays can be achieved.\nThe organic EL element emits light when current flows through an organic emissive layer. In many cases, a hole transport layer and an electron transport layer formed of an organic material are provided adjacent to the organic emissive layer so as to assist light emission. These organic layers are apt to degrade when exposed to moisture.\nAccordingly, in organic EL displays, in order to remove moisture from the organic layers, a cathode formed of a metal is provided for covering the upper portion of the organic EL element while a space above the display region (a region where pixels are provided) in which the organic EL element is disposed is made air proof and a drying agent is provided in this space.\nHowever, there remains a demand for more reliable prevention of diffusion of moisture into the organic layers of the organic EL element. Although it is possible to increase the thickness of the cathode as a solution, for example, when the thickness is too large, the cathode would be deformed and peeled due to the resulting stress."}
{"text": "iBiquity Corporation has developed a specification for its “in-band on-channel” (IBOC®) broadcasting system that meets the requirements of the Federal Communications Commission (FCC). Transmitting a hybrid (both analog and digital) IBOC®-compatible broadcast requires radiating an analog signal with frequency modulation (FM) technology and a digital signal with orthogonal frequency division multiplexing (OFDM) technology. The OFDM signal occupies the edges of the FM signal's emissions mask and has a total radiated power one hundredth (−20 dB) that of the FM signal. Each hybrid IBOC® signal uses one of the hundred radiotelephone channels for public reception established between television channels 6 and 7 in the very high frequency (VHF) band (88.1 MHz to 107.9 MHz). IBOC® also defines standards for all-digital VHF and for AM-band (535 KHz to 1705 KHz) radio.\nA previous IBOC® antenna design disclosed in U.S. Pat. No. 7,084,822 (“the '822 patent”), incorporated herein by reference, includes crossed dipoles for radiation of analog and digital signals. The propagation concept disclosed includes, in at least one embodiment, two pairs of dipoles in each bay, with the dipoles in each pair spaced horizontally by a quarter wavelength, oriented at right angles to each other within parallel planes, and driven with two substantially unrelated signals, where the two signals are fed as traveling waves from opposite ends of a coaxial line and coupled therefrom to drive the dipoles.\nA crossed-dipole pair so driven reinforces signal emission at some azimuths and cancels signal emission at other azimuths to produce generally peanut-shaped and overlaid circularly polarized patterns—beams—for the two signals. Each beam has two lobes; the lobes for that beam have the same circular polarization, but are opposite in phase at each instant. The '822 patent discloses a second dipole pair that taps the coaxial line a quarter wavelength from a first dipole pair for impedance cancellation, and that has an azimuthal orientation at right angles to that of the first pair, so that each bay radiates two circularly polarized signals with opposite handedness and oppositely rotating phase. The signals generally fill in at intermediate azimuths to an extent sufficient for the antenna to be termed omnidirectional.\nWhile effective, this embodiment is somewhat constrained by the traveling-wave feed method, and is better suited to tower-top mounting and a small number of bays. A second embodiment in the '822 patent feeds crossed dipole pairs from taps on a traveling wave coaxial line, splitting the tapped signals to drive the pairs. This allows all of the radiating elements to be placed to one side of the coaxial line, but is still further limited in power by halving the number of coupling taps per radiator.\nAnother previous IBOC® antenna design is disclosed in copending U.S. application Ser. No. 11/698,065, filed Jan. 26, 2007, titled “Antenna System and Method to Transmit Cross-Polarized Signals from a Common Radiator with Low Mutual Coupling,” incorporated herein by reference. This design includes separate corporate feed from analog and digital transmitters to a plurality of hybrid couplers per bay, each hybrid including unbalanced inputs and balanced outputs, so that multiple crossed-dipole radiators with integral cross-coupling cancellation can be provided in a plurality of bays with low mutual coupling. While highly effective, broad banded (>20% BW for VSWR<1.05:1), and high power capable, this design can be complex, preferably using either a tower-top mounting scheme or a plurality of discrete mountings around a tower or other structure to realize omnidirectional coverage.\nMultiple-channel broadcast towers are costly to build and occupy significant amounts of real estate in rare locations (high up and near the center of population regions but low in local population, so transmitters can be clustered around them). Many such broadcast towers are relatively full, that is, they are limited in the number of antennas that can be mounted on them with adequate vertical separation, and desirable positions such as tower tops are typically already taken, leaving small or low positions or replacement of existing antennas as enhancement possibilities. Some IBOC®-compatible antenna designs are not readily adapted to tower-side mounting, because they use highly symmetrical structures to achieve omnidirectional patterns and would require robust, extended—and massive—cantilever brackets for tower side mounting."}
{"text": "As set out in the above referenced applications/patents, the Applicant has spent a substantial amount of time and effort in developing printheads that incorporate micro electro-mechanical system (MEMS)—based components to achieve the ejection of ink necessary for printing.\nAs a result of the Applicant's research and development, the Applicant has been able to develop printheads having one or more printhead chips that together incorporate up to 84 000 nozzle arrangements. The Applicant has also developed suitable processor technology that is capable of controlling operation of such printheads. In particular, the processor technology and the printheads are capable of cooperating to generate resolutions of 1600 dpi and higher in some cases. Examples of suitable processor technology are provided in the above referenced patent applications/patents.\nCommon to most of the printhead chips that the Applicant has developed is a component that moves with respect to a substrate to eject ink from a nozzle chamber. This component can be in the form of an ink-ejecting member that is displaceable in a nozzle chamber to eject the ink from the nozzle chamber.\nA particular difficulty that the Applicant has been faced with is to achieve a suitable interface between a prime mover in the form of an actuator and the moving component. This interface is required to permit the moving component to be displaced in the nozzle chamber and to inhibit leakage of ink from the nozzle chamber.\nAs set out in the above referenced patents/patent applications, the printhead chip is manufactured using integrated circuit fabrication techniques. This is the usual manner in which MEMS-based devices are fabricated. Such forms of fabrication are subject to constraints since they involve successive deposition and etching techniques. It follows that MEMS-based devices are usually formed in layers and that components having relatively complex shapes are difficult and expensive to fabricate.\nIn FIG. 1, reference numeral 10 generally indicates part of a nozzle arrangement of a printhead chip. The part 10 shown illustrates an actuator 12 and an ink-ejecting member 14. The actuator 12 includes an elongate actuator arm 16 that extends from an anchor 18. The actuator arm 16 is configured so that, when it receives a drive signal, the actuator arm 16 bends towards a substrate 20 as indicated by an arrow 22. A connecting formation 24 is interposed between the actuator arm 16 and the ink-ejecting member 14. Thus, when the actuator arm 16 is bent towards the substrate 20, the ink-ejecting member 14 is displaced in the direction of an arrow 26 to eject ink from the nozzle chamber.\nIt would be intuitive simply to use the arrangement 10 together with a suitable sealing structure to achieve effective ink ejection and sealing. The reason for this is that it would appear that the actuator arm 16, the connecting formation 24 and the ink-ejecting member 14 could be in the form of a unitary structure. However, the Applicant has found that it is not possible to achieve a working configuration as shown by using MEMS-based fabrication techniques. In particular, it has been found by the Applicant that such a unitary structure does not lend itself to such fabrication techniques.\nIt follows that the Applicant has been led to conceive the present invention."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a Metal Oxide Semiconductor Field Effect Transistors (MOSFETs) made on semiconductor on insulator (SOI) substrates. In particular, the invention relates to MOSFETs with a recessed channel film in an SOI layer. The recessed channel film forms an abrupt junction.\n2. Description of Related Art\nU.S. Pat. No. 7,041,538 B2, issued on May 9, 2006 to Ieong et al., describes a high performance CMOS device on an SOI substrate with a gate recessed into an SOI layer and ion implanted source/drain regions with halo and extension implants.\nU.S. Pat. No. 6,939,751 B2 issued on Sep. 6, 2005 to Zhu et al. describes a raised source drain field effect devices with a channel recessed into a silicon germanium film located above a SOI layer.\nU.S. Pat. No. 7,652,332 B2 issued on Jan. 26, 2010 to Cartier et al. describes extremely thin silicon on insulator transistor with raised source/drain, high dielectric constant (high-k) oxide and metal gate.\nU.S. Pat. No. 7,429,769 B2 issued on Sep. 30, 2008 to Diaz et al. describes a recessed channel field effect transistor (FET).\nIn a paper entitled “Extremely Thin SOI (ETSOI) CMOS with Record Low Variability for Low Power System-on-Chip Applications,” by K. Cheng et al. published in 2009 IEEE International Electron Device Meeting, Dec. 7-9, 2009 a method of making CMOS transistors on ETSOI substrates is disclosed.\nIn a paper entitled “Ultra-thin SOI replacement gate CMOS with ALD TaN/high-k gate stack,” by B. Doris et al. in IEEE VLSI-TSA International Symposium on VLSI Technology, Apr. 25-27, 2005 a device built on an ETSOI substrate using a replacement gate substrate is disclosed.\nIn a presentation by B. Doris et al., entitled “FD SOI for Low Power CMOS,” 2009, available at http://www.soiconsortium.org/pdf/fullydepletedsoi/FD%20SOI%20for%20Low%20Power%20CMOS.pdf, a summary of device performance challenges and possible solutions are reviewed. Some possible solutions include various devices made using ETSOI substrates.\nIn a paper entitled “CMOS Transitions to 22 and 15 nm,” by D. Lammers published Jan. 1, 2010, in Semiconductor International describes device structures and possible methods of manufacture for FETs at ground rules less than or equal to 22 nm. Possible devices include planar MOSFETs on ETSOI substrates.\nIn a paper entitled “Tri-Gate Transistor Architecture with High-k Gate Dielectrics, Metal Gates and Strain Engineering,” by J. Kavalieros et al. available at http://download.intel.com/technology/silicon/tri-gate_paper_VLSI—0606.pdf, the authors describe a non-planar MOSFET on SOI with recessed source and drains."}
{"text": "Many factors make hemodialysis inefficient, difficult, and expensive. These factors include the complexity of hemodialysis, the safety concerns related to hemodialysis, and the very large amount of dialysate needed for hemodialysis. Moreover, hemodialysis is typically performed in a dialysis center requiring skilled technicians. Therefore any increase in the ease and efficiency of the dialysis process could have an impact on treatment cost or patient outcome.\nFIG. 1 is a schematic representation of a hemodialysis system. The system 5 includes two flow paths, a blood flow path 10 and a dialysate flow path 20. Blood is drawn from a patient. A blood flow pump 13 causes the blood to flow around blood flow path 10, drawing the blood from the patient, causing the blood to pass through the dialyzer 14, and returning the blood to the patient. Optionally, the blood may pass through other components, such as a filter and/or an air trap 19, before returning to the patient. In addition, in some cases, anticoagulant may be supplied from an anticoagulant supply 11 via an anticoagulant valve 12.\nA dialysate pump 15 draws dialysate from a dialysate supply 16 and causes the dialysate to pass through the dialyzer 14, after which the dialysate can pass through a waste valve 18 and/or return to the dialysate feed via dialysate pump 15. A dialysate valve 17 controls the flow of dialysate from the dialysate supply 16. The dialyzer is constructed such that the blood from the blood flow circuit flows through tiny tubes and the dialysate solution circulates around the outside of the tubes. Therapy is achieved by the passing of waste molecules (e.g., urea, creatinine, etc.) and water from the blood through the walls of the tubes and into the dialysate solution. At the end of treatment, the dialysate solution is discarded."}
{"text": "The invention relates to a process for the study of chemical reactions in the presence of potentially catalytic substances, in which the reactions are carried out in parallel in reactors. The invention further relates to an apparatus, in particular for carrying out this process, the apparatus having a plurality of parallel-connected reactors which are provided with feedlines and discharge lines.\nProcesses and apparatuses of said type are known and are used, inter alia, in the search for catalysts for the heterogeneous or homogeneous catalysis of industrial chemical processes.\nHowever, most recently, novel techniques have made it posible to prepare substances in a great number, which could be potential catalysts for a multiplicity of chemical processes (P. G. Schultz et al., Science 1995, 1738). The study of this multiplicity of potential catalysts is scarcely possible any longer using the conventional series screening processes, since these screening processes are limited with respect to throughput and analytical resolution, as well as in reproducibility. Frequently, for pure activity screening, integral effects such as heating of the catalyst etc. are used completely inadequately without direct product mixture analysis or efficiency analysis. Furthermore, optimizing the conditions of catalyst activation and process procedure make special requirements of quantitative analytical processes and of the reproducibility of the reaction conditions.\nThe object underlying the invention was therefore to develop an inexpensive process or an inexpensive apparatus by which a multiplicity of chemical reactions can be studied in a short time and thus reproducible, qualitative and quantitative data can be obtained with respect to the composition of the different reaction mixtures and reaction products.\nThis object is achieved by a process of the type mentioned at the outset, which is characterized in that the reactions are carried out in miniaturized reactors and the reaction mixture or the reaction products are analyzed with respect to nature and amount during the reaction period.\nThe object is further achieved by an apparatus of said type, which is characterized in that the reactors are miniaturized, at a volume from the range from 0.001 cm3 to 1 cm3.\nThe invention thus relates to a process for the study of chemical reactions in the presence of potentially catalytic substances, in which the reactions are carried out in parallel in reactors, characterized in that the reactions are carried out in miniaturized reactors and the reaction mixture is analyzed with respect to nature and amount during the reaction period.\nThe invention further relates to an apparatus, in particular for carrying out this process, the apparatus having a plurality of parallel-connected reactors which are provided with feedlines and discharge lines, characterized in that the reactors are miniaturized, at a volume from the range from 0.001 cm3 to 1 cm3.\nParticular embodiments or developments of the invention result from the corresponding subclaims. Individual or a plurality of individual features mentioned in the claims can also each themselves represent solutions of the invention, and the features within the claim categories can also be combined as desired.\nA particular embodiment of the process according to the invention is characterized in that the reaction starting materials are fed continuously to the reactors and the reaction products are discharged continuously from the reactors. However, the batch mode of operation is equally possible.\nA further particular embodiment is characterized in that use is made of starting materials which are at least in part labeled with isotopes, preferably with deuterium (2H) or heavy oxygen (18O) or heavy carbon (13C) or mixtures thereof. These generate characteristic spectral shifts in rotation-vibration spectra, which, in addition to labeling the reaction pathway, owing to starting material mixed variants, can lead to novel interesting reactions or reaction products and relatively small by-product contents can be contrasted systematically.\nIndividual reactors or reactors combined to form groups can also be fed with different starting material mixtures, in order to reveal or detect any synergies present in this manner using the methods of combinatory chemistry.\nThe starting material mixtures, reaction mixtures or product mixtures can be analyzed with respect to nature and amount of the substances present at any time points of the course of the reaction by spectroscopic analysis, preferably by infrared spectroscopy (IR), particularly preferably by Fourier IR spectroscopy. Other spectroscopic methods, such as laser spectroscopy or UV spectroscopy, are equally suitable for the study. The process can be carried out at differing temperatures and pressures, at temperatures of the range from xe2x88x9250xc2x0 C. up to and including 600xc2x0 C., preferably from room temperature to 500xc2x0 C., or at various pressures, at absolute pressures from 10xe2x88x923 to 103 bar, preferably from 10xe2x88x922 to 200 bar. The data obtained can then be supplied to a comprehensive parametric and data analysis.\nThe invention is further characterized in that the reactions can be carried out in the presence of a heterogeneous or homogeneous catalyst and in that it is possible to screen the catalytic activity (ie. product detection) and selectivity (main product distribution) of catalyst amounts less than 10 mg, preferably less than 1 mg, in a reactor.\nIn a particular embodiment of the apparatus according to the invention, a plurality of miniaturized reactors separated from one another can be disposed in a block. The volume of these reactors can be in the range from 0.001 cm3 to 1 cm3, preferably from 0.01 cm3 to 0.5 cm3, particularly preferably from 0.05 cm3 to 0.2 cm3. In a further preferred embodiment of the apparatus according to the invention, the reactors are disposed as a square or rectangular pattern in a metal block, which can be parallelepipedal or cubic. The metal block can be provided with heating block elements or cooling elements and can be fitted with a temperature sensor in the vicinity of each reactor. This makes monitored and reproducible temperature control possible. For example, by this means, a defined temperature gradient can be set over the metal block. The reactors are advantageously disposed in a plane which is parallel to a surface of the parallelepiped. Advantageously, at least some of the feedlines and discharge lines of the individual reactors are perpendicular to this plane. They can be designed as through-holes in the metal block. The reactors can be designed as boreholes. The number of reactors in one block can be greater than 20, preferably greater than 40, particularly preferably greater than 100, very particularly preferably greater than 200. Using these reactors, under defined reaction conditions, a small amount of potential catalysts (also termed samples below) can be brought into contact and reacted in parallel, ie. simultaneously, with starting material or starting material mixtures in liquid and/or gaseous form batchwise or continuously. Automation of the apparatus according to the invention is possible, in particular charging the reactors with catalysts can proceed automatically, preferably by a laboratory robot or a pipetting machine.\nIn a further preferred embodiment, the miniaturized reactors are designed in the metal block as 4 mm bore holes and disposed in such a manner that differing starting material and inert gases can flow through them through 2.5 mm capillary bore holes. The gases then pass into a spacer, preferably a distancing plate, which is mounted on the metal block and in which the bore holes of the metal block continue. The arrangement of metal block and spacer is provided with a customary cuvette bore hole in which the gases can be analyzed by spectroscopy. For this purpose, the bore hole is sealed at both ends with a transparent window. If infrared spectroscopy is to be used for analysis, use is preferably made of windows of 1-1-1 silicon, NaCl, KBr, Ge, ZnSe or KSR5. For the analysis, a collimated analytical beam, in the case of IR spectroscopy an infrared beam, is decoupled from an interferometer, preferably without a mirror, and directed through a dry-gas-flushed space through the cuvette bore hole onto a detector situated behind. The cuvette bore hole can be, for example, 5 mm thick. By choosing a suitable thick spacer, the length of the cuvette bore hole can be chosen between a few cm (1-10) and several 10 s of cm (10-50), depending on the reaction conditions and reaction type. To record the spectra, the analytical beam can be directed through all cuvette bore holes successively by a deflection apparatus. However, a plurality of beams or a plurality of analyzers can also be used, so that simultaneous recording of spectra is possible with a plurality of reactors. However, equally, the block containing the reactors can be moved by motion devices, for example stepper motors, in such a manner that all cuvette bore holes are brought successively into the beam path of the spectrometer. Suitable materials for block and spacer are preferably the customary corrosion-resistant metallic materials familiar to those skilled in the art, particularly aluminum or steel, preferably stable to rust and/or acid and/or high temperature.\nA further embodiment of the apparatus according to the invention which is particularly suitable for homogeneous catalysis is characterized in that, in the case of at least one reactor having a volume which is preferably less than 200 xcexcl, an ATR crystal (attenuated total reflection spectroscopy, preferably conically pointed, preferably made of ZnSe or KSR5 or diamond) makes possible spectroscopic contact with the reaction mixture for differing solvents and reaction conditions and pressures up to 200 bar. In this case, the analytical beam is focused onto the ATR crystal.\nThe advantages of the process according to the invention and the apparatus according to the invention are essentially that a multiplicity of reactions can be carried out rapidly, inexpensively and reproducibly under virtually identical conditions and with comparatively small amounts of substance and sample and, in the course of this, can simultaneously be studied spectroscopically. It thus offers the opportunity of using the opportunities discussed in connection with combinatorial chemistry (K. Burgess et al., Ang. Chem. 1996, 108, 2, 192, incorporated by reference into the application) for industrial catalyst screening. By choosing identical samples and differing other reaction conditions such as temperature, pressure, starting material composition, a parallel reaction optimization can be carded out.\nTo analyze the data obtained, advantageously, a data matrix is prepared in such a manner that all selectable and documentable reaction conditions (starting material partial pressures, starting material composition, temperature, flow or flow rate, total pressure, sample composition, sample lattice parameters and all base points of the spectra) are given according to reaction conditions, ie. for each reactor, as columns of the matrix. This matrix can be subjected to a factor analysis (E. R. Malinowski et al., Factor Analysis in Chemistry, Wiley, New York, 1980, incorporated by reference into the application), by calculating the covariance matrix, the eigenvalues, the abstract eigenvectors, the loadings and the coefficients of the multidimensional regression, and preferably outputing them as files. A (pre)normalization of the data can also be selected by the mean xe2x80x9c0xe2x80x9d and standard deviations xe2x80x9c1xe2x80x9d, which can avoid base line effects or absolute effects. This permits the prediction of differing variables from calibration data sets (eg. quantitative CO2 contents at different temperatures), determination of the dependence of parameters in spectral ranges to optimization of the analysis, generation of different distance matrices from the output data (for example the similarity of catalysts with respect to selected variables and properties) and direct feedback of the catalyst composition to a laboratory robot synthesizer which mixes a set of new catalyst samples and synthesizes them xe2x80x9cautomaticallyxe2x80x9d by sintering or calcining on a robot line.\nAn embodiment of the process according to the invention and a development of the apparatus according to the invention are described in more detail below with reference to FIGS. 1 and 2, without the intention of restricting the invention in any manner."}
{"text": "Interrupting structures which have been utilized in conventional starter or contractor equipment, such as suitable for 400 amperes, 5 KV interrupting equipment, are well known in the art. Generally, such equipment is utilized in connection with fuse structures, and under certain (severe) conditions it is desirable to maintain the contactor contacts closed, so that the fuse equipment will fuse or \"blow\" first. It is usual in such starter equipment to provide a quick-opening and quick-closing with adequate contact pressure existing in the closed-circuit position of the equipment.\nIt is also desirable, in such starter equipment, to provide optional equipment which will enable the customer to specify either an air-break-type of interrupter or a vacuum-type interrupter. However, the latter type of equipment involves particular requirements, with which the present invention is concerned. U.S. Pat. No. 3,264,433 issued Aug. 2, 1966 to Russell D. Clark, Jr. and Charles J. Mahlar entitled \"Switching Apparatus Employing Fuse As Movable Contacts and Having Safety Features\" involves equipment related to the subject matter of the present invention. Also U.S. Pat. No. 3,264,431, issued Aug. 2, 1966, to Alfred W. Hodgson, entitled \"Compact Switching Apparatus Employing Fuse As Movable Contacts\", and assigned to the assignee of the instant invention relates to similar type equipment. Also U.S. Pat. No. 3,290,468, issued Dec. 6, 1966 to Russell D. Clark, Jr., and Charles J. Mahlar, entitled \"Clip For Fuse Employed As Movable Member Switch\" relates to the present equipment.\nU.S. Pat. No. 3,264,431, issued Aug. 2, 1966 to Alfred W. Hodgson is concerned with the equipment of the type utilized in practicing the invention. Finally, reference may be had to U.S. Pat. No. 3,264,432, issued Aug. 2, 1966 to Alfred W. Hodgson and Russell D. Clark, Jr., as well as U.S. Pat. Nos. 3,898,407 and 3,921,109 to Alfred W. Hodgson, which related to similar-type equipment, all of the foregoing patents being assigned to the assignee of the instant invention, and all of the aforesaid patents relating to air-break equipment."}
{"text": "The field of the present invention is the field of Leading Zero Anticipators (LZA) that predict the number of leading zeros or leading ones in a sum of mantissas generated by an adding device used to add two floating point numbers.\nLeading Zero Anticipators (“LZAs”) are typically used to predict the number of leading zeros in a sum of mantissas during floating point addition. Prediction of the number of leading zeroes is used to normalize the result of the addition and is performed at the same time as (in parallel with) mantissa addition, which increases the speed of the normalization process that is performed after mantissa addition has been completed.\nIn a traditional LZA design shown in FIG. 1, mantissaA and mantissaB (labeled as 10) both enter the adder 13 and the LZA 12 simultaneously. A normalization shift is performed by the shifter 15 after addition is performed. LZAs 12 are commonly used in devices performing fast addition that try to reduce the number and length of clock cycles required to do the addition and in devices that implement a single instruction which performs multiplication followed by addition. The use of redundant formats such as the carry-save format during addition does not require the use of the LZA during the addition. But the use of redundant formats does require placement of an LZA in parallel with a second adder which adds the sum and carry bit streams generated by the carry-save addition process. The sum and carry bit streams are added before the final normalization stage is completed. In such circuits, the carry and sum become the input operands (mantissaA and mantissaB) that are transmitted to the adder in FIG. 1. The nature of addition in redundant formats invalidates the use of several LZAs in the final add operation. For instance, in the carry-save format, the carry term is shifted left by 1 bit position (equivalent to multiplying the carry value by 2) before addition is performed. If there is a negative carry term and a positive sum term, the left shift that is performed after carry save addition might shift out the leading one of the carry term, making it a positive number (the leading bit indicates sign). The addition of two seemingly positive numbers (sum and carry) may then produce a valid negative number. This seemingly contradictory result is valid in carry-save addition but invalidates the assumptions made in the designs of many LZAs. For example, consider a carry (value −6) and a sum (value 7) with expected final result −6×2+7=−5. The corresponding 4 bit binary vectors in 2's complement are carry 1010 and sum 0111. When the carry is shifted left 1 bit, it becomes 0100. Adding 0100 (which is now seemingly a positive number) to 0111 (also a positive number) gives 1011 which is a negative number (−5) and is the correct expected final result.\nAnother drawback of conventional LZAs occurs if the operation is a subtraction. Traditional application of LZAs during subtraction requires that the mantissa of the smaller operand must be subtracted from the mantissa of the larger operand so that a positive result is generated. To insure that the operand of smaller magnitude is subtracted from the greater one, a comparator is used to determine which mantissa is larger. This is expensive because of the need to use extra logic gates. Another traditional method of leading zero anticipation does not require comparison but uses 2 LZAs together which are implemented so that they make opposite assumptions about the sign of the generated result when anticipating the number of leading zeroes or ones. The correct LZA output is then chosen depending on the sign of the result of the addition. This is also expensive.\nYet another conventional implementation of LZAs does not require comparison of mantissas and uses just one LZA. However logic designers implementing this form of LZA must assume that the “normal” rules of addition apply when the operands are added. Consequently, this implementation cannot be used when the input operands are expressed in redundant formats, such as the carry save format. This is the case because addition in redundant format does not follow all the conventions of normal addition as explained above.\nTherefore, what is required is a novel LZA that may be implemented with a single LZA and that is implemented in a way that does not make assumptions regarding rules of addition that preclude the use of redundant formats such as the carry save format. Also, an LZA that does not require comparison of the magnitude of the input operands or knowledge of the sign of the result in order to anticipate the number of leading zeroes or leading ones is needed."}
{"text": "Nowadays, with advancements in technologies such as virtual reality, augmented reality and so forth, users of specialized devices are able to experience digitally simulated environments (namely, virtual worlds). Specifically, such simulated environments enhance a user's perception of reality around him/her by providing the user with a feeling of complete involvement in the simulated environments using contemporary techniques such as stereoscopy. Moreover, such simulated environments relate to fully virtual environments (namely, virtual reality) as well as real-world environments including virtual objects therein (namely, augmented reality).\nTypically, the user may utilize specialized devices, such as an augmented reality device, for experiencing such simulated environments. For example, the augmented reality device may be a binocular augmented reality device (for example, such as augmented reality glasses) having one display per eye, and cameras to capture a real-world environment of a user. In other words, the displays of the binocular augmented reality device show different two-dimensional images of virtual objects to the left and right eyes of the user for making the virtual objects appear realistic from a perspective of the user. Examples of such augmented reality devices include near-eye displays such as augmented reality headsets, augmented reality lenses, and so forth.\nHowever, conventional augmented reality devices have certain limitations. For example, the size of displays suitable for closely imitating visual acuity of the human eyes is too large to be accommodated within the conventionally available augmented reality devices that aim to be moderately small. Specifically, displays with field of view comparable to the human eyes are dimensionally very large, and therefore, not suitable for use in such augmented reality devices.\nMoreover, presently available augmented reality devices are unable to achieve optimum trade-off between the field of view and angular resolution, thereby, compromising on either the field of view, or on the angular resolution. Furthermore, the presently available augmented reality devices, namely see-through augmented reality devices, are operable to render the simulated environments in a limited field of view and/or a limited angular resolution. Consequently, the presently available augmented reality devices are limited in their capability to render realistic and/or highly detailed virtual objects in a large field of view.\nFurthermore, conventional augmented reality devices are unable to resolve vergence-accommodation conflict, thereby causing strain in the eyes of the user. The cause is typically two-fold: 1. Inability to cover convergence of eye needed to focus at close-by objects (as displays are too small to cover enough horizontal field of view to cater for both focus at long and short distance) and 2. Inability to change optical focus distance to reflect the actual real world distance. Therefore, the conventional augmented reality devices are limited in their ability to mimic the human visual system.\nTherefore, in light of the foregoing discussion, there exists a need to overcome the aforementioned drawbacks associated with conventional augmented reality devices."}
{"text": "1. Technical Field\nThe entire disclosure of Japanese Patent Application No.: 2013-039456, filed Feb. 28, 2013 is expressly incorporated by reference herein.\nThe present invention relates to a technique that is related to cartridges for storing a liquid.\n2. Related Art\nThere is a conventionally known technique of employing an ink cartridge for storing ink (also simply called a “cartridge”) as a technique for supplying ink to a printer, which is one example of a liquid ejection device. This cartridge includes a liquid storage portion for storing ink serving as the liquid, and a liquid supply portion for supplying the ink in the liquid storage portion to a printer.\nThere is also known to be a cartridge that includes a circuit substrate for exchanging information with a printer (e.g., see JP-A-2008-074100). This circuit substrate stores information regarding the cartridge, such as information indicating the manufacturing date of the cartridge and the color of stored ink. When the cartridge is mounted to the cartridge mounting portion of the printer, terminals of the circuit substrate come into contact with conductive contact members arranged on the cartridge mounting portion. The circuit substrate and the printer are thus electrically connected, making it possible for various types of information to be exchanged between the circuit substrate and the printer.\nWhen the cartridge is in the mounted state, it is subjected to various external forces. One example of an external force is force that the conductive contact members exert on the cartridge (biasing force). Another example of an external force is the force exerted on the cartridge when the printer performs a printing operation or the like. If the printer is of the type in which the cartridge mounting portion is installed on the carriage that is provided with the head (i.e., the “on-carriage type”), another example of an external force is the force exerted on the cartridge due to the cartridge moving along the main scanning direction of the carriage (inertial force).\nWhen the cartridge is subjected to various external forces, the state of contact between the contact members and the terminals becomes unstable, and there are cases where the electrical connection between the cartridge and the printer cannot be maintained in a favorable state.\nIn the technique described in JP-A-2008-074100, a positioning portion is provided on the bottom plane of the cartridge in order to suppress positional shift caused by force applied to the cartridge by the conductive contact members. The positioning portion of the cartridge is shaped as a projection, and engaging it with a positioning portion provided on the cartridge mounting portion suppresses positional shift of the cartridge caused by force applied by the contact members.\nHowever, with the technique described in JP-A-2008-074100, there have been cases where it is difficult to suppress positional shift of the terminals provided on the cartridge relative to the contact members caused by external force applied to the cartridge. For example, there have been cases where the reliability of contact between the contact members and the terminals decreases due to the terminals and the positioning portions being arranged at separated positions. In this way, in technology that employs cartridges provided with terminals, there is desire for a technique for raising the reliability of contact between the terminals provided on the cartridge and the contact members provided in the printer. Also, with cartridges provided with terminals, with liquid ejection devices, and with systems that include a cartridge and a liquid ejection device, there is desire for size reduction, cost reduction, resource saving, structure simplification, improvement in usability, lifetime extension, improvement in safety during use, and the like."}
{"text": "One of the fundamental problems which exist in long haul high rate optical communications is chromatic dispersion of light pulses passing through fiber optic lines. This dispersion causes different wavelengths to travel through the optical waveguide at different speeds. An optical communications light pulse which is part of a bit stream is created by a transmission laser operating at a predetermined carrier optical frequency or wavelength (the words frequency, represented by f, and wavelength in vacuum, represented by l, will be used interchangeably, as is currently done in the art, with l related to f by the equation l=c/f, where c is the speed of light, viz. 299 792 458 m/s).\nThe pulse has a given intensity and duration which can be detected at a receiver located at an opposite end of the communications waveguide (typically the optic fiber) some distance away from the transmitter, perhaps a few to tens of kilometers away. The light pulse comprises spectral components covering a certain bandwidth Df, which is generally in the range 10-60 GHz in modern optical communications systems operating in the bit rate range of 2.5-10 Gb/s. In optical fibers used for communications, the group velocity associated with each spectral component varies with its optical frequency (or wavelength). This means that the lowfrequency spectral components of a light pulse do not travel at the same velocity as its high frequency components. This causes a light pulse initially injected into a fiber with duration Dt to broaden out (in others words to become \"dispersed\") to a duration Dt' after propagating through a certain length of optical fiber. As an example, in conventional fibers deployed over land, the dispersion is typically 17 ps/nm/km. This means that two spectral components separated by say 10 GHz at a carrier frequency of 193 000 GHz (l=1553.329 nm), corresponding therefore to a wavelength separation of 0.0804 nm, would undergo a differential delay (or dispersion) of 17 ps/nm/km times 0.0804 nm, i.e. 1.37 ps/km, or 137 ps after 100 km of fiber (the sign of the dispersion in conventional fibers is such that low frequency light travels at a lower group velocity than highfrequency light). Light pulses that may be initially 25 ps in duration in a 10 Gb/s optical communication system have spectral components that cover several tens of GHz. The time dispersion between the low- and the high-frequency spectral components of such a light pulse are therefore more than 137 ps after a 100-km stretch of fiber. Such a large amount amount of dispersion is intolerable since, the pulses being 100 ps apart, dispersion will make subsequent pulses start overlapping to a considerable degree and increase the bit error rate.\nTherefore, dispersion causes a spreading of laser pulses over great distances (even in the most non-dispersive optic fibers) such that the pulse energy is spread out and lower-frequency spectral components trail behind higher-frequency components of the laser pulse.\nWhen the bit time slot in an optical communication system is large enough, dispersion causes a small portion of the pulse energy to spill over into adjacent bit time slots. The pulse height at the receiver will be marginally lower, but still detectable as a bit, and empty bit slots (e.g. representing zeros) will contain small amounts of spill over light, which will be below the predetermined detection threshold. However, as the bit time slots are made smaller so as to increase the bit rate and increase transmission capacity, dispersion of a very short light pulse may significantly reduce the light intensity in the bit time slot and even cause enough spill over into adjacent empty bit time slots so as to cause difficulty in the detection of zeroes.\nDispersion thus produces the spreading of short optical pulses, and thus affects negatively the quality of a communication link. If the dispersion of a light pulse is great enough, the communication link utilizing the light pulses becomes unstable, and eventually, unusable.\nIn addition to problems in optical communications, dispersion is also a major problem in the generation of very short, high power optical laser pulses. The spreading of the pulses reduces the achieved peak power, and thus reduces the efficiency of the laser pulses.\nVarious attempts have been made to address the problem of dispersion in optical signals. Three basic approaches have been developed. The first approach is to pass the optical signal transmitted, which has suffered dispersion due to the optical transmission waveguide, through a length of waveguide having an opposite sign of dispersion at the communications wavelength before feeding the signal to the receiver. U.S. Pat. No. 4,969,710 to Tick et al. is an example of such a method. The second approach has been to use optical devices to separate the light pulse into its wavelength components, then to suject the separated components to different delays before recombining the components into a single dispersion compensated optic signal. U.S. Pat. No. 5,473,719 to Stone is an example of this approach. The third approach has been to use chirped in-fiber Bragg gratings to reflect each wavelength component at different points such that the reflected optical signal has been dispersion compensated. This approach is disclosed in a paper by Franqois Ouellette, titled \"Dispersion cancellation using linearly chirped Bragg grating filters in optical waveguides\", Optics Letters, vol. 12, pp. 847-849, October 1987. Experimental results have been reported in the paper by W. H. Loh, R. I. Laming, N. Robinson, A. Cavaciuti, F. Vaninetti, C. J. Anderson, M. N. Zervas, and M. J. Cole, titled \"Dispersion compensation over distances in excess of 500 km for 10-Gb/s systems using chirped fiber gratings\", IEEE Photonics Technology Letters, vol. 8, pp.944-946, July 1996.\nOther known prior art attempts to resolve these problems arm disclosed in U.S. Pat. No. 5,570,439 to Ido et al.; U.S. Pat. No. 5,568,583 to Akasaka et al; U.S. Pat. No. 5,559,920 to Chraplyvy et al.; U.S. Pat. No. 5,530,778 to Rast; U.S. Pat. No. 5,504,829 to Evans et al.; U.S. Pat. No. 5,448,674 to Vengsarkar et al.\nKnown dispersion compensators are bulky, generally limited to a single channel wavelength and/or have significant losses. The advantage of the dispersion management device proposed here is that it could me manufactured using integrated optic techniques so that it could be very compact and low-cost. In addition the taper lends itself very naturally to tayloring the dispersion as a function of wavelength in a way to match (and compensate) that of an optical fiber. In 10 Gb/s and higher bit rate system the variation of the dispersion with wavelength (from 16 ps/nm/km at 1530 nm to 17 ps/nm/km at 1580 nm, for example, in certain optical fibers) is a factor that is now considered. Our device can offer this wavelength taylored dispersion profile and can do so over very broad wavelength bands."}
{"text": "Passwords may be used in many situations where a user needs to be authenticated. Situations requiring authentication may include the ability to access or use a computer, mobile phone, PDA, or any other device; they may also enable access to a physical location, or allow use of a credit/debit card or similar instrument. Passwords are typically alpha-numeric strings or sequences entered on a keyboard.\nLike reference numerals are used in the drawings to denote like elements and features."}
{"text": "Mobile communication devices, especially mobile phones, have become a convenient means for information exchange. Recently, their function has been extended from a mere communicational tool to an information collection tool as well. One example is the mobile phone with a built-in digital camera. The pictures taken by the camera can be transmitted in real time via a wireless communication system and shared among a group of people.\nMore functions can be built into a mobile communication device by utilizing its information transmission capability.\nIn U.S. Pat. No. 6,947,571 to Rhoads et al (hereinafter merely Rhoads), describes a cell phone equipped with a 2D optical sensor, which enables a variety of applications. For example, such a phone may also be provided with a digital watermark decoder, permitting decoding of steganographic data on imaged objects. Movement of a phone may be inferred by sensing movement of an imaged pattern across the optical sensor's field of view, allowing use of the phone as a gestural input device through which a user can signal instructions to a computer-based process. A variety of other arrangements by which electronic devices can interact with the physical world are also detailed, e.g., involving sensing and responding to digital watermarks, bar codes, RFIDs, etc.\nIn U.S. Patent Application No. 2003/0187742 to Junichi Yamagishi, (hereinafter merely Yamagishi), describes a program for personal authentication using a mobile phone that enables the mobile phone user to buy a commodity without cash and a card, and without entering a password. A personal authentication system includes a mobile phone, a reading adapter and an authentication server. The mobile phone can download a program for personal authentication from a web server. The mobile phone can, according to the downloaded program, output a train of serial optical pulses that represents personal information. The reading adapter can receive and read the output pulse train. The authentication server registers the personal information and performs personal authentication based on the read pulse train. In particular, Yamagishi describes a personal authentication system, in which a mobile phone is employed to download a personal authentication program and output a train of optical pulses according to the program to represent personal information.\nIn U.S. Patent Application No. 2003/0136837 to Amon et al (hereinafter merely Amon) describes a method and a system for the local or remote authentication of an item, in particular a security document, with the help of a authenticating device, comprised in, connected to, or linked to mobile communication equipment. Said item carries a marking exhibiting a characteristic physical behavior in response to interrogating energy, such as electromagnetic radiation and/or electric or magnetic fields. Said marking may comprise physical and logical security elements, e.g. a barcode, or a characteristic particle or flake pattern, exhibiting a characteristic physical response. In other words, a method is disclosed to use a mobile communication device for item authentication. The items, in particular security documents are authenticated through the methods of imaging, spectroscopy, barcode, radio-frequency resonance, etc. As can be seen, among various embodiments of the Amon patent application, only the imaging method utilizes a camera that integrated in the mobile phone. In other embodiments, the authentication devices are accessories that are linked to the mobile phone through wired or wireless connections.\nIn U.S. Patent Application No. 2005/0111723 to Hannigan et al, digital watermarking methods and systems are described. In one embodiment, a handheld device displays a digitally watermarked image at an event center. An optical scanner captures optical scan data representing the device display. A watermark decoder analyzes the captured optical scan data to obtain the watermark. The watermark is used to determine authorized entry. The handheld device may include a cell phone or other portable computing device. In other words, a handheld device, such as a cell phone, is used to display a digitally watermarked image for authorized entry.\nIn U.S. Patent Application No. 2006/0082760 to Lin describes an optical sensing module, an optical sensing and image capturing architecture, and a method for optically scanning fingerprints with a portable communications system associated with a mobile phone to capture a group of images of a fingerprint on a finger of the user of the mobile phone. The optical sensing module has a light permeable device, at least one light source, and a sensing device. The sensing device has a plurality of sensors that form a line array of sensors. The optical sensing and image capturing architecture has the above optical sensing module and an image output unit. The volume of the optical sensing module is reduced for facilitating installation in the portable communications system, for fragmentally capturing a group of images of the fingerprint, and for processing of identification of the fingerprint. In other words, an optical sensing and image capturing architecture is applied to a mobile phone to capture a group of images of a fingerprint for personal identification.\nNone of the above mentioned references disclosed a method of integrating an optical spectroscopic sensor with a mobile communication device and the numerous applications engendered thereof, which will be the main topic of the present invention."}
{"text": "In modern equipment used for information processing, rapid printers in which printing of the characters is accomplished without causing relief-printing type to strike a sheet of receptive paper are being used more and more. These printing machines, known as non-impact printers, generally include a recording medium, most frequently comprising a rotating drum or an endless belt, which is provided with a recording surface on which sensitized zones corresponding to characters to be printed can be created by electrostatic or magnetic means, these zones being capable of attracting the solid particles of a powdered developer product.\nVarious applicator devices can be used for applying the solid developer particles to the recording carrier of a printing machine of this type. For example, the device which is described and shown in U.S. Pat. No. 3,161,544 may be utilized. However, despite all the care taken in realizing this device, it is difficult to prevent the developer particles from being deposited not only to excess on the sensitized zones of the recording medium but also, although in very slight quantity, outside these zones. This phenomenon has been attributed to the fact that if the particles are charged with humidity or static electricity, or if they undergo melting, which, however slight, makes them more or less sticky, they adhere to the surface with which they have been placed in contact. An excessive deposit of developer particles on the sensitized zones of the recording medium is undesirable because when this developer is transferred to the sheet of receptive paper, there is a risk that the developer that has been deposited in accordance with the configuration of the image formed by these sensitized zones may smudge and blur the image. On the other hand, the deposit of particles outside the sensitized zones of the recording medium is also undesirable, because when these particles are transferred to the paper they form a base which reduces the contrast between the transferred image and the original base of the paper.\nIn the prior art, various retouching devices have been used to eliminate excess developer on the surface of the recording medium. One example of a device of this kind has been particularly described in U.S. Pat. No. 3,816,799, in which a mass of developer particles, placed in contact with the surface of the recording medium and downstream of a device for applying the particles, discharges the particles which adhere to the recording medium outside the sensitized zones, such that these discharged particles become detached from the recording medium and clump together with this mass. In any event a retouching device of this kind is not completely satisfactory in use, because it does not always assure a complete electrical discharge of the particles and as a consequence does not enable the reliable elimination of the developer particles which remain on the recording medium outside the sensitized zones.\nAnother known retouching device is described in U.S. Pat. No. 3,643,629 and includes a magnetic assembly rotating in the interior of a fixed cylindrical sleeve. The assembly itself is embodied by a plurality of magnetic elements, in the form of a sector, disposed side by side about a rotating shaft. These elements are magnetized in such a way as to have peripheral magnetized zones on their peripheral surface, the magnetic polarity of which remains constant along a direction parallel to the shaft, but alternating in passing from one peripheral zone to the next. Nevertheless, with this retouching device in which the magnetic element accommodated in the interior of the fixed sleeve rotated past the recording medium in the opposite direction from that in which the medium is displaced, uniform removal of the excess particles is not obtained, because in certain parts of the medium the removal is too pronounced and in other parts it is insufficient.\nThis disadvantage can be overcome by using the retouching device described in U.S. Pat. No. 4,067,018, which enables the elimination of the excess developer particles on an endless magnetic tape on which a latent magnetic image has been produced. This retouching device includes a guide roller of radius R, over which the endless magnetic tape travels. When this tape is driven at a constant speed V, the developer particles located on this tape are subjected as they travel past the guide roller to the action of a centrifugal force, the value of which, for each particle of mass M, is equal to MV.sup.2 /R. Under these conditions, the intensity of this centrifugal force can be adjusted by suitably selecting the speed at which the tape is driven and the radius of the guide roller such that the value of the centrifugal force is greater than that of the forces retaining the particles on the zones at the base of the tape, yet is not so great as to be able to detach all the particles located on the magnetized zones of this tape. The excess developer particles which have thus been detached from the tape are then eliminated by means of an aspirator disposed near the guide roller.\nIt has been found that to remove the excess developer particles on this magnetic tape the acceleration imparted by this centrifugal force must be at least equal to 10 g, g being the value of the acceleration of gravity. As a result, in the case where the guide roller used in this retouching device has a diameter of twenty millimeters, the endless magnetic tape must be driven at a speed of at least one meter per second. A displacement speed on this order, however, means that this retouching device cannot be used in modern magnetic printing machines, in which the drive speed of the recording medium in practice does not exceed 30 cm/s. To enable the use of this retouching device in a printing machine of this kind, the diameter of the guide roller must accordingly be reduced, but to do so means that a guide roller must be selected the diameter of which must not exceed two millimeters. A guide roller of that size would necessarily be fragile, and its use in a magnetic printing machine would not offer all the security that is desired."}
{"text": "1. Field of the Invention\nThis invention relates generally to valve arrangements for beverage dispensers in which a diluent (e.g., carbonated water) is mixed with a beverage syrup, and more specifically, this invention relates to a quick disconnect arrangement for such beverage dispenser valves which prevents removal of the dispenser valve for cleaning, repair, or replacement while the syrup and diluent is being supplied to the valve under pressure.\n2. Description of the Prior Art\nBeverage dispensers are in standard use in restaurants and in vending machines to mix and provide soft drink beverages such as colas and the like upon demand. Such soft drinks typically comprise a beverage syrup which includes all of the ingredients and flavorings mixed with a diluent such as carbonated water. Beverage dispensers are preferred over premixed bottled or canned beverages in restaurants and in some vending machines because of the greatly reduced storage requirements for unmixed syrup as opposed to the finished canned or bottled beverage.\nHowever, because such beverage dispenser valves are electromechanical devices usually operated by solenoids such as the fast flow valve arrangement illustrated in U.S. Pat. No. 4,266,726--Brown, et al., and because the beverage syrup contains numerous ingredients including a high concentration of sugar, it is often necessary to remove the valve arrangement for cleaning and/or repair.\nFurther, since the carbonated water and the beverage syrup are supplied to the dispensing valve under pressure, it is desirable to provide a disconnect arrangement which permits easy, rapid removal of the dispensing valve but which assures that the dispensing valve cannot be removed while the carbonated water and beverage syrup are being supplied to the dispensing valve under pressure. Very often in restaurants and other fast food establishments, multiple valve arrangements for different beverages are being supplied by the same source of carbonated water. Thus, it is desirable to provide a disconnect arrangement which does not require a shutdown of the entire system to permit removal of a single valve arrangement."}
{"text": "When an electronic document such as a webpage is displayed, a part of the image may be magnified.\nThis application provides an information-processing device includes a display controller that, when a scaling operation for changing a scale of an image to be displayed is not performed, displays the image with a first magnification in a display, and when the scaling operation is performed, displays the image with a second magnification that is greater than the first magnification in the display."}
{"text": "The electric power grid has become more distributed. More localized electrical generation facilities operated by users produce an ever greater share of electricity. As the prevalence of localized electrical general facilities increases, more and more electronic loads localized throughout the grid increase problems of power factor (PF).\nThe PF typically refers to the relationship between current and voltage in an alternating current (AC) electric power system. The PF of an AC electric power system can be defined as the ratio of the real power flowing to a load to the apparent power that is seemingly being drawn from a power source. PF typically is a number between 0 and 1, and is frequently expressed as a percentage. Real power can be understood as the capacity of a circuit to perform work in a particular time. Apparent power can be understood as the product of the current and the voltage of the circuit."}
{"text": "In electrophotography, also known as Xerography, electrophotographic imaging or electrostatographic imaging, the surface of an electrophotographic plate, dram, belt or the like (imaging member or photoreceptor) containing a photoconductive insulating layer on a conductive layer is first uniformly electrostatically charged. The imaging member is then exposed to a pattern of activating electromagnetic radiation, such as light. The radiation selectively dissipates the charge on the illuminated areas of the photoconductive insulating layer while leaving behind an electrostatic latent image on the non-illuminated areas. This electrostatic latent image may then be developed to form a visible image by depositing finely divided electroscopic marking particles on the surface of the photoconductive insulating layer. The resulting visible image may then be transferred from the imaging member directly or indirectly (such as by a transfer or other member) to a print substrate, such as transparency or paper. The imaging process may be repeated many times with reusable imaging members.\nAn electrophotographic imaging member may be provided in a number of forms. For Example, the imaging member may be a homogeneous layer of a single material such as vitreous selenium or it may be a composite layer containing a photoconductor and other materials. In addition, the imaging member may be layered in which each layer making up the member performs a certain function. Current layered organic imaging members generally have at least a substrate layer and two electro active or photo active layers. These active layers generally include (1) a charge generating layer containing a light-absorbing material, and (2) a charge transport layer containing hole transport molecules or materials. These layers can be in a variety of orders to make up a functional device, and sometimes can be combined in a single or mixed layer. The substrate layer may be formed from a conductive material. Alternatively, a conductive layer can be formed on a nonconductive inert substrate by a technique such as but not limited to sputter coating.\nThe charge generating layer is capable of photogenerating a charge and injecting the photogenerated charge into the charge transport layer or other layer.\nIn the charge transport layer, the hole transport molecules may be in a polymer binder. In this case, the hole transport molecules provide hole or electron transport properties, while the electrically inactive polymer binder provides mechanical properties. Alternatively, the charge transport layer can be made from a charge transporting polymer such as vinyl polymer, polysilylene or polyether carbonate, wherein the charge transport properties are chemically incorporated into the mechanically robust polymer.\nImaging members may also include a charge blocking layer(s) and/or an adhesive layer(s) between the charge generating layer and the conductive substrate layer. In addition, imaging members may contain protective overcoatings. These protective overcoatings can be either electroactive or inactive, where electroactive overcoatings are generally preferred. Further, imaging members may include layers to provide special functions such as incoherent reflection of laser light, dot patterns and/or pictorial imaging or subbing layers to provide chemical sealing and/or a smooth coating surface.\nImaging members are generally exposed to repetitive electrophotographic cycling, which subjects the exposed charge transport layer or alternative top layer thereof to mechanical abrasion, chemical attack and heat. This repetitive cycling leads to a gradual deterioration in the mechanical and electrical characteristics of the exposed charge transport layer.\nAlthough excellent toner images may be obtained with multilayered belt or drum photoreceptors, it has been found that as more advanced, higher speed electrophotographic copiers, duplicators and printers are developed, there is a greater demand on copy quality. A delicate balance in charging image and bias potentials, and characteristics of the toner and/or developer, must be maintained. This places additional constraints on the quality of photoreceptor manufacturing, and thus, on the manufacturing yield.\nDespite the various approaches that have been taken for forming imaging members, there remains a need for improved imaging member design, to provide improved imaging performance, longer lifetime, and the like."}
{"text": "The present invention relates generally to the field of disposable articles utilized for absorption and containment of urine and other body exudates. More particularly, the invention relates to a disposable article, for example a diaper, that incorporates components such as an absorbent insert pad which may be removed from the article and flushed down a toilet.\nDisposable articles for the absorption and containment of urine and other body exudates are generally known in the art. Such disposable articles have found particular utility in the fields of infant care, child care, feminine care, and adult incontinency. Present commercially available disposable articles for such uses are generally unitary, preshaped or prefolded, and comprised of a fluid pervious bodyside liner, a fluid impervious backing sheet, and an absorbent material disposed between the bodyside liner and the backing sheet. Disposable articles of this type effectively and efficiently absorb and contain urine or other body exudates. Such articles are designed for a single use and are simply deposited after use in a waste receptacle. This ease and convenience of disposal has contributed to the increasing popularity of such disposable absorbent articles.\nUntil now, the vast majority of the absorbent articles of :he aforementioned type have been designed to be disposed of by delivery to solid waste landfills. Such absorbent articles have not generally been compatible with being flushed down a toilet. The relatively few products that have been designed specifically to be disposed of, at least partially, by flushing have been ineffective or impractical.\nWith regard to diapers, for example, U.S. Pat. No. 3,211,147 to P. O. Pherson et al. discloses a diaper pad that is intended to be completely disposed of in household toilets and septic tanks. The pad includes a paper web, several fluff pads, and paper reinforcing strips. The diaper employs numerous channels to distribute fluids. Because the pad contains only wood pulp derivatives, the user may separate the diaper into pieces and flush the pieces in a toilet.\nA diaper with a flushable portion is disclosed in U.S. Pat. No. 3,667,466 to H. J. Ralph. This patent discloses a diaper having a flushable layer and a reusable retainer. The flushable layer is composed of top and bottom layers of wet-strength material that surround an absorbent layer. The diaper incorporates numerous perforations which cause the diaper to be self-segmenting and self-disintegrating when immersed in a toilet.\nAnother diaper with a flushable portion is disclosed in U.S. Pat. No. 4,964,857 to C. Osborn. This patent discloses a diaper that incorporates a removable and flushable inner sheet of material. The inner sheet has a liquid permeable layer, an absorbent layer, and a moisture repellent layer. The inner sheet, which is attached to the remainder of the diaper by adhesives, is constructed with a plurality of perforations and a series of hollow channels. A drawstring located within the channels is used to remove the inner sheet from the remainder of the diaper.\nThese attempts to provide diapers that are completely or partially flushable have not been entirely successful. For example, the above-mentioned completely flushable diapers must necessarily be limited to materials that can be flushed. Consequently, completely flushable diapers have provided less than desirable fit and absorbency characteristics. On the other hand, the diapers having a flushable layer and a reusable retainer are frequently messy and inconvenient. Finally, the diapers that incorporate a removable insert have not effectively utilized the absorbent material contained within the insert and have tended to prematurely experience saturation spots. Consequently, these diapers have not provided an adequate amount of comfort and dryness. Additionally, these diapers have included devices such as drawstrings to remove the inserts. Devices of this type present numerous manufacturing difficulties and increase the expense and complexity of production.\nThe present invention teaches an improved disposable article which provides components, such as an absorbent insert pad, which may be quickly and conveniently removed from the article. The user can dispose of the insert pad, as well as any solid feces material held in the article, in any toilet that is connected to a holding tank, a municipal sewer, or other municipal wastewater treatment system. The insert pad may then be flushed down the toilet, and the remainder of the article may be composted or recycled where suitably equipped facilities are available, disposed at a conventional landfill site, or incinerated. Alternately, the complete absorbent article including the insert pad may be directed into the solid waste stream. Thus, the user may select among several alternative disposal methods, thereby accommodating the individual's preferences and the community's capabilities."}
{"text": "This invention relates to apparatus for determining the opening combination of a safe lock and, in particular, to a programmable device capable of automatically dialing an ordered sequence of combinations into the lock and testing the lock after each combination has been set in to determine if the combination will open the lock.\nWhere the combination of a safe or vault has been forgotten, lost or otherwise unobtainable, it is generally left to the skill of the individual locksmith who is employed to open the safe to determine exactly what digits are embodied in the opening combination. Even under ideal conditions, where the workings of the lock are familiar and one or more numbers in the combination are known, the task of opening the safe can be extremely time consuming and costly. In some procedures, particularly where the workings of the lock mechanism are old and not understood, holes are sometimes drilled into the lock assembly to allow the locksmith to actually see and/or manipulate the tumblers and release mechanism embodied therein. These drilling procedures while effective in opening the safe or vault usually make it necessary to replace the lock. In any event, regardless of the procedures involved, finding the opening combination of a safe or vault has typically proven to be an extremely difficult and sometimes frustrating task."}
{"text": "1. Field of the Invention\nThe present invention relates to a method, system, program, and data structures for testing a network system including input/output (I/O) devices.\n2. Description of the Related Art\nA storage area network (SAN) comprises a network linking one or more servers to one or more storage systems. Each storage system could comprise a Redundant Array of Independent Disks (RAID) array, tape backup, tape library, CD-ROM library, or JBOD (Just a Bunch of Disks) components. Storage area networks (SAN) typically use the Fibre Channel Arbitrated Loop (FC-AL) protocol, which uses optical fibers to connect devices and provide high bandwidth communication between the devices. In Fibre Channel terms the switch connecting the devices is called a “fabric”. The link is the two unidirectional fibers, which may comprise an optical wire, transmitting to opposite directions with their associated transmitter and receiver. Each fibre is attached to a transmitter of a port at one end and a receiver of another port at the other end. When a fabric is present in the configuration, the fibre may attach to a node port (N_Port) and to a port of the Fabric (F_Port).\nBecause a Fibre Channel storage area network (SAN) is an amalgamation of numerous hosts, workstations, and storage devices, troubleshooting for errors can often be a somewhat complex process. Currently, in the prior art, a technician will perform a series of tests from a host system in the SAN and test various channels and connections to the storage devices to detect problems and then try to locate the specific source of a problem. Technicians generally rely on their own knowledge, experience and expertise when diagnosing the SAN system for errors. Such knowledge is not a shared resource, but rather an individual point of view and an accumulation of guess work and personal experience. As a result, it is unlikely that different storage experts troubleshoot a storage system in the same manner, thereby leading to possible incorrect or inconsistent diagnosis as well as an increase in the Mean Time To Diagnose (MTTD). Moreover, as the number of SAN systems proliferate, it may become more and more difficult for system administrators to locate available diagnosticians.\nCertain “cookbook” approaches to testing a Fibre Channel network have been proposed, such as the “Fibre Channel FC-AL-2 Parametric Test Suite Rev. 7.0”, published by the Fibre Channel Consortium, document no. ANSI X3.272-199X (January 2000), which publication is incorporated herein by reference in its entirety. Such documents describe specific tests that may be performed to troubleshoot a Fibre Channel network. However, again the order in which the tests are selected and performed is still a matter of choice for the diagnostician performing the troubleshooting operations.\nNotwithstanding current efforts at troubleshooting network components, such as a SAN, the current art lacks tools that provide an integrated and consistent approach toward diagnostic testing of a SAN and its components."}
{"text": "1. Field of the Invention\nThis invention generally relates to optical devices and a method for making such devices. More particularly this invention is related to devices for viewing the interior of a passageway in a selective region.\n2. Description of Related Art\nThere are a variety of applications which require an individual to view the interior of a remote chamber. In an industrial environment, for example, it is desirable to view portions of the interior of a machine, such as a jet engine, that is otherwise inaccessible. In medical applications it is desirable to view a passageway in the body. Viewing scopes that enable one to view such remote passages are well known in the art. Such devices used for mechanical applications are generally called borescopes. For medical applications the viewing scopes take many forms and are known by different names. Generically the name \"endoscope\" means any slender, tubular optical instrument used as a viewing system for examining an inner part of the body. In practice such viewing scopes have a variety of names, such as endoscopes, laparoscopes and bronchoscopes. In the following discussion the phrase \"viewing scope\" is meant to include generically both (1) borescopes and (2) endoscopes in any form including laparoscopes, bronchoscopes, etc.\nViewing scopes generally have an eyepiece at a proximal end, an objective lens at a distal end and an image guide interposed between the eyepiece and the objective lens. The image guide uses two general forms of optics, namely: rigid optics and fiber optics. Image guides using rigid optics include relay or rod lenses to transfer an image from a distal end of the viewing scope to its proximal end. The following patents disclose examples of image guides using such rigid optics:\nU.S. Pat. No. 4,036,218 (1977) Yamashita et al PA0 U.S. Pat. No. 4,993,817 (1991) Hoogland PA0 U.S. Pat. No. 4,776,668 (1988) Fujimoto PA0 U.S. Pat. No. 4,938,205 (1990) Nudelman PA0 U.S. Pat. No. 4,945,894 (1990) Kawashima PA0 U.S. Pat. No. 4,961,738 (1990) Mackin\nImage guides using fiber optics generally interpose a bundle of optical fibers between the proximal and distal ends of the viewing scope. The following patents disclose image guides that incorporate such fiber optics:\nThe Fujimoto patent discloses a viewing scope with a coherent fiber optic bundle as an image guide. Each fiber comprises a pure quartz silica strand coated with a reflective layer. This patent discloses structure eliminating differential expansion between the fiber optic bundle and other portions of the viewing scope during flexure.\nThe Nudelman patent discloses an endoscope that includes a fiber optic bundle between distal and proximal ends. The fiber optic bundle is flexible between the ends.\nThe Kawashima patent discloses an endoscope in which image guide fibers transmit an image formed at the objective to a proximal end of the endoscope.\nThe Mackin patent discloses an angioplasty catheter that includes a fiber optic bundle for conveying an image between the distal and proximal ends.\nConventionally fiber optic bundles used in such image guides comprise a preferred three-glass fiber or two-glass fiber. A three-glass fiber has a glass core that acts as a light transmitter, a concentric reflective layer and a fusible, acid soluble outer layer. In conventional processing, a bundle of such fibers are assembled together and heated to fuse just the outer layer. Then the ends are isolated and the entire structure is immersed in an acid bath as a leaching solution to remove portions of the fusible glass between the ends so the bundle will be flexible along its length. It is difficult to control this leaching process. If the leaching continues for too long a time, the solution acid can etch through the outer glass layer and the reflective layer thereby detracting from the image by lowering its contrast. This leaching step effectively controls the minimum diameter of the individual fibers and the number of individual fibers that can be formed into a bundle of a given size. As the diameter of the fibers becomes smaller, the effort required to control the leaching process increases because the fusible layer becomes thinner. Consequently costs increase because the manufacturing is more difficult and because production yields decrease.\nOptical fibers used in two-glass fiber optic bundles have a light transmitting core and a reflective outer cladding. Multiple drawing steps can reduce the light transmitting core to under 10 microns. During each drawing step a bundle of fiber passes through a die that reduces both its core and cladding and that causes the bundles and fibers within each bundle to adhere. When such fibers are drawn below 10 microns, however, the cladding becomes very thin and transmissive. Cross talk can occur and limit the contrast resolution of the viewed image. In some situations this requires the insertion of black glass or other barriers to prevent light from one fiber from transferring to another fiber through the thin cladding.\nThe basic criteria for selecting rigid optics or fiber optics for an image guide include spatial resolution as a primary criteria and contrast resolution, cost and ruggedness. Rigid optics provide the best spatial and contrast resolution. However, they are fragile and have the highest manufacturing costs. An image guide must be designed separately for each combination of a viewing scope length and diameter. Fiber optics, on the other hand, provide a rugged image guide and require less design effort and have lower costs. However, image guides using fiber optic bundles have limited spatial resolution. In viewing scopes of the type contemplated by this invention, the spatial resolution is limited to 5,000 to 30,000 pixels in an image. Fiber optics also provide somewhat lower contrast resolution than rigid optics and require an inventory of one bundle for each combination of viewing scope diameter and length."}
{"text": "The present disclosure relates generally to internet protocol television (IPTV), and more particularly, to methods, systems, and computer program products for providing IPTV support.\nCurrently, IPTV services are provided to users through IPTV networks and in-home devices, such as set top boxes (STBs). When a user experiences difficulty in receiving the IPTV services, the existing process is for the user to call a customer service representative to seek assistance. The customer service representative can often respond to routine issues, but may have difficulty responding to more complex issues. This often results in the IPTV service provider sending a technician to the user's location to diagnose and resolve the issue. Unfortunately, the task of sending technicians to a user's location is costly and time consuming for both the service provider and the user.\nThere is a need in the art for methods and systems to facilitate diagnosis and resolution of IPTV service issues without requiring sending a technician to visit the user."}
{"text": "Current networks typically include devices, such as routers, switches or gateways, which transfer or switch data, such as packets, from one or more sources to one or more destinations. A packet is one format of data in which encapsulated data can be transmitted through a network based on control information that is stored in a header portion of the packet. A router is a switching device that can receive a packet and, based on the packet header, may forward the packet towards its final destination.\nExisting routers include forwarding engines for receiving and forwarding incoming packets to their intended destinations. To forward incoming packets from an input port to an appropriate output port, routers may perform complex data manipulation actions on the packet header. Such data manipulation actions frequently result in the router rewriting a portion of the packet header before transmitting the packet to an appropriate output port of the router."}
{"text": "Conventionally, for the purpose of preventing the occurrence of operating noises of a compressor inside a room, a compressor unit has been provided outside the room and helium gas discharged from the compressor has been air-cooled by an air-cooling heat exchanger. However, since the cooling by the air-cooling heat exchanger installed outside the room is air-cooling, the helium gas cannot be cooled to a temperature lower than the air temperature by this air-cooling heat exchanger. Therefore, when the outdoor temperature is high in summer or in a similar case, the temperature of the helium gas supplied to the cryogenic expander is hard to be maintained at a temperature not higher than a temperature (35.degree. C., for example) at which the insulating property of a motor is assured.\nIn view of the above, there is proposed a cryogenic refrigerating apparatus (Japanese Patent Laid-Open Publication No. HEI 6-249148) which can cool the helium gas to a lower temperature through two-stage cooling by installing a compressor unit provided with a compressor and a first air-cooling heat exchanger outside a room and installing an intermediate unit provided with a second air-cooling heat exchanger inside the room.\nIn this cryogenic refrigerating apparatus which performs the two-stage cooling, as shown in FIG. 7, a compressor unit 1 is constructed of a helium gas compressor 11, a first air-cooling heat exchanger 12 comprised of, for example, a cross fin coil provided in communication with discharge side piping 21 of this compressor 11 and an oil separator 13 provided in communication with the discharge side piping 21 on the outlet side of this first air-cooling heat exchanger 12, while an intermediate unit 3 provided with a second air-cooling heat exchanger 31 comprised of, for example, a cross fin coil is provided separately from this compressor unit 1, the compressor unit 1 being installed outside a room and the intermediate unit 3 being installed inside the room.\nAn end portion which belongs to the discharge side piping 21 connected to the discharge side of the compressor 11 is connected to gas supply piping 41 of the intermediate unit 3, while an end portion which belongs to an intake side piping 22 connected to the intake side of the compressor 11 is connected to gas return piping 42 of the intermediate unit 3.\nThe gas supply piping 41 of the intermediate unit 3 is connected to a high-pressure side communication piping 51 communicated with a cryogenic expander 5, while the gas return piping 42 of the intermediate unit 3 is connected to a low-pressure side communication piping 52 communicated with the cryogenic expander 5.\nFurther, the second air-cooling heat exchanger 31 is connected to the gas supply piping 41, an adsorber 32 is provided in communication with the outlet side of the second air-cooling heat exchanger 31 and the second air-cooling heat exchanger 31 is provided with a fan 33.\nIt is to be noted that oil collected in a bottom portion of the oil separator 13 is injected into a compression element of the compressor 11 via oil injection piping 23 and oil collected to a height higher than a specified oil surface height inside the oil separator 13 is returned from the intake side piping 22 into the compressor 11 via an oil return piping 24. On the other hand, oil collected in a bottom portion inside the compressor 11 is cooled in the first air-cooling heat exchanger 12 via an oil cooling piping 25 and thereafter returned from the intake side piping 22 into the compressor 11.\nThen, in the first air-cooling heat exchanger 12 of the compressor unit 1 installed outside the room, by making compressed high-temperature helium gas to exchange heat with the outdoor air to firstly cool it by the outdoor air for the achievement of the greater part of the heat radiation of the helium gas outside the room and further cooling the helium gas by the second air-cooling heat exchanger 31 of the intermediate unit 3 installed inside the room, the compressed helium gas is cooled in two steps by the outdoor air and the indoor air. This arrangement has allowed the helium gas to be cooled to a temperature not higher than a specified temperature (35.degree. C., for example) even when the outdoor temperature is high and prevented the operating noises inside the room with the compressor unit 1 installed outside the room.\nIn the cryogenic refrigerating apparatus constructed as above, the helium gas, of which cooling has been insufficient in the compressor unit 1, can be cooled in the intermediate unit 3. However, in regard to the cooling in the intermediate unit 3, constant cooling is consistently performed no matter whether the cooling capacity of the first air-cooling heat exchanger 12 of the compressor unit 1 depending on the outside air temperature is great or small, i.e., a fan 33 for cooling the second air-cooling heat exchanger 31 is driven to rotate consistently at a constant rotating speed so as to make the air flow constant. Therefore, when the outdoor temperature is low in winter, constant cooling is performed by the second air-cooling heat exchanger 31 regardless of the load from outside the room in spite of the fact that sufficient cooling has been performed in the first air-cooling heat exchanger 12, and this has resulted in excessive cooling and a significant change in refrigerating capacity, causing a disadvantage that a stable refrigerating operation can still not be performed.\nFurthermore, although not shown in FIG. 7, the first air-cooling heat exchanger 12 of the compressor unit 1 is normally provided with an outdoor fan for cooling. When starting the refrigerating apparatus in a case where the outdoor temperature is extremely low in winter, the viscosity of the oil (mainly ether-based oil) inside the units 1 and 3 is very high. Therefore, when excessive cooling is performed by the operation of the outdoor fan at the first air-cooling heat exchanger 12, the viscosity of the oil does not reduce, and this has tended to cause a disadvantage that the units 1 and 3 do not correctly operate.\nThe present invention is intended to solve the aforementioned problems, and its principal object is to provide a cryogenic refrigerating apparatus capable of reducing the change in refrigerating capacity as far as possible with respect to a wide range of change in outdoor temperature and performing a stable refrigerating operation.\nAnother object is to allow the units to regularly operate by speedily reducing the viscosity of the oil even when the outdoor temperature becomes very low in winter."}
{"text": "Many different enterprises run complex networks of servers to implement various automated communication functions to the enterprise. For example, as mobile wireless communications have become increasingly popular, carriers such as Verizon Wireless™ have customer communication systems to provide notifications of account related activities to their customers, for example, as Short Message Service (“SMS”) messages to account holders' mobile stations, as emails, etc. Because of the large number of customers served by a major carrier, and the level of account activities, the volume of notification message traffic is quite large. To effectively provide such notifications, Verizon Wireless™ implemented its Customer Communication Enterprise Services (“CCES”) as an enterprise middleware web service.\nAt a high level, the CCES middleware comprises a web server layer and an application layer. The architecture allows clients to send a request, for example for a notification, to a web server. The http web server then forwards the client request to one of a number of application servers. Each application server has multiple applications running on it. The application server determines the proper application to process the client request based on the context root of the client request. The application server processes the client request, in the CCES example, by sending one or more request messages to a back end system such as the Vision Billing System, Mobile Telephone Activation System (“MTAS”), the SMS gateway and others, for example, to implement account activity and to initiate subsequent automatic notification thereof to the account holder. Once the application server has processed the request, a reply is then sent back to the web server which will then forward the reply back to the client.\nIn a multiple server environment described above, there are many servers running many applications that perform a variety of functions. These applications may require constant updating and modifications. The updates may be done on scheduled basis once a week, for example. The process of updating these applications with new code/configuration may generally be the same from week to week with subtle but important differences. This process is called a “Change Control” at, for example, Verizon Wireless™. At Verizon Wireless™ there are about seven groups of servers, which may run the same applications, for implementing the CCES functionality. Each group of server may have anywhere from eight to sixteen physical machines. There are four steps to the current software “Change Control” process.\nThe first step in installing new or updated server software includes installing updated application code. In this step, an updated version of the application code is to be deployed on the servers in the group. This is involves: (a) taking down the affected servers running the file to be updated, (b) backing up current server application code, (c) installing the new An Enterprise ARchive (“EAR”) file, (d) restarting the server (e.g., each machine has its own GUI to do this process), and (e) notifying the testers to begin once the server application is back on line. The EAR is a file format used for packaging one or more modules into a single archive so that the deployment of the various modules onto an application server happens simultaneously and coherently.\nThe second step includes installing updated application configuration files. In about 70% of the update deployments several files may need to be deployed/updated. These are supporting application configuration files or simply files that customers may download depending on the application being used by the customer. The process to accomplish this is by copying files to some/all of the machines in the group.\nThe third step includes executing Structured Query Language (“SQL”) files. After performing the first and second steps, the SQL may need to be executed. This may be done once and may not be per machine type process. The deployer may be given some SQL files and may run them against three different Oracle tables. These may be run on one central Graphical User Interface (“GUI”).\nThe fourth step includes running of scripts. In some of the server groups, scripts may need to be run at various points within the deployment process. These scripts are small programs that need to be run at specific points/steps within the process to perform functions that aid in the deployment.\nThe steps in the above process are manual and require the deployer to log onto individual machines in the server group to perform the steps sequentially. That is, deploying two to three groups of servers would involve manually logging into twenty-four machines and running multiple process steps on each. This requires more than one person because one person cannot keep up with all the tasks that are occurring simultaneously. Therefore, there is a need to improve on the time-consuming repetitive tasks requiring two or more human resources to perform the four step deployment process described above, for example, on as many as twenty-four machines."}
{"text": "Engine torque may be increased in response to a change in a requested or desired engine torque. Engine torque may be increased via increasing an amount of air and fuel supplied to the engine, at least up to a condition where the engine lacks capacity to induct additional air into engine cylinders. The amount of fuel supplied to the engine may be increased as the amount of air inducted to the engine increases such that the engine operates near stoichiometric conditions. However, there may be conditions where the amount of air inducted to engine cylinders is less than desired. For example, if the engine is supplied gaseous fuel, the gaseous fuel may displace some air in the cylinder, thereby limiting the amount of torque the engine may produce. Consequently, a desired engine torque response may not be provided.\nThe inventor herein has recognized the above-mentioned disadvantages and has developed an engine operating method, comprising: directly injecting a gaseous fuel to an engine in response to a pressure of gaseous fuel stored in a fuel tank exceeding a threshold level; and increasing an idle speed of the engine and injecting the gaseous fuel to an engine air intake in response to the amount of gaseous fuel stored in the fuel tank being less than the threshold level.\nBy increasing idle speed of an engine in response to a pressure of fuel stored in a fuel tank, it may be possible to improve an engine's torque response. For example, if a gaseous fuel is directly injected to engine cylinders, the engine may respond quickly to an increase in desired engine torque. However, if the engine transitions to injecting the gaseous fuel via engine intake injection (e.g., injecting fuel to an engine intake manifold or intake ports) when higher pressure gaseous fuel is not available from the gaseous fuel storage tank, engine output torque at engine idle may not respond as quickly as is desired. The engine torque response at lower engine speeds may be dampened since the engine may generate less torque at lower engine speeds because a portion of cylinder volume may be displaced by fuel. However, if engine idle speed is increased, engine power output may increase more rapidly in a short period of time since more combustion events occur for each second of engine operation. Thus, performance of an engine having engine intake injection may be improved so that vehicle acceleration may be more consistent between direct injection and engine intake injection modes.\nThe present description may provide several advantages. In particular, the approach may reduce improve engine torque response. Further, the approach may make changes between fuel injection modes less noticeable to a driver. Further still, the approach may extend a vehicle's travel range while operating on a gaseous fuel.\nThe above advantages and other advantages, and features of the present description will be readily apparent from the following Detailed Description when taken alone or in connection with the accompanying drawings.\nIt should be understood that the summary above is provided to introduce in simplified form a selection of concepts that are further described in the detailed description. It is not meant to identify key or essential features of the claimed subject matter, the scope of which is defined uniquely by the claims that follow the detailed description. Furthermore, the claimed subject matter is not limited to implementations that solve any disadvantages noted above or in any part of this disclosure."}
{"text": "1. Field of the Invention\nThe present invention relates to a golf club head with a face component composed of a metal material, and an aft-body composed of a light-weight material. More specifically, the present invention relates to a golf club head with face component composed of a metal material for a more efficient transfer of energy to a golf ball at impact, and a non-metallic aft-body to control the mass distribution.\n2. Description of the Related Art\nWhen a golf club head strikes a golf ball, large impacts are produced that load the club head face and the golf ball. Most of the energy is transferred from the head to the golf ball, however, some energy is lost as a result of the collision. The golf ball is typically composed of polymer cover materials (such as ionomers) surrounding a rubber-like core. These softer polymer materials having damping (loss) properties that are strain and strain rate dependent which are on the order of 10-100 times larger than the damping properties of a metallic club face. Thus, during impact most of the energy is lost as a result of the high stresses and deformations of the golf ball (0.001 to 0.20 inch), as opposed to the small deformations of the metallic club face (0.025 to 0.050 inch). A more efficient energy transfer from the club head to the golf ball could lead to greater flight distances of the golf ball.\nThe generally accepted approach has been to increase the stiffness of the club head face to reduce metal or club head deformations. However, this leads to greater deformations in the golf ball, and thus increases in the energy transfer problem.\nSome have recognized the problem and disclosed possible solutions. An example is Campau, U.S. Pat. No. 4,398,965, for a Method Of Making Iron Golf Clubs With Flexible Impact Surface, which discloses a club having a flexible and resilient face plate with a slot to allow for the flexing of the face plate. The face plate of Campau is composed of a ferrous material, such as stainless steel, and has a thickness in the range of 0.1 inches to 0.125 inches.\nAnother example is Eggiman, U.S. Pat. No. 5,863,261, for a Golf Club Head With Elastically Deforming Face And Back Plates, which discloses the use of a plurality of plates that act in concert to create a spring-like effect on a golf ball during impact. A fluid is disposed between at least two of the plates to act as a viscous coupler.\nYet another example is Jepson et al, U.S. Pat. No. 3,937,474, for a golf Club With A Polyurethane Insert. Jepson discloses that the polyurethane insert has a hardness between 40 and 75 shore D.\nStill another example is Inamori, U.S. Pat. No. 3,975,023, for a Golf Club Head With Ceramic Face Plate, which discloses using a face plate composed of a ceramic material having a high energy transfer coefficient, although ceramics are usually harder materials. Chen et al., U.S. Pat. No. 5,743,813 for a Golf Club Head, discloses using multiple layers in the face to absorb the shock of the golf ball. One of the materials is a non-metal material.\nLu, U.S. Pat. No. 5,499,814, for a Hollow Club Head With Deflecting Insert Face Plate, discloses a reinforcing element composed of a plastic or aluminum alloy that allows for minor deflecting of the face plate which has a thickness ranging from 0.01 to 0.30 inches for a variety of materials including stainless steel, titanium, KEVLAR®, and the like. Yet another Campau invention, U.S. Pat. No. 3,989,248, for a Golf Club Having Insert Capable Of Elastic Flexing, discloses a wood club composed of wood with a metal insert.\nAlthough not intended for flexing of the face plate, Viste, U.S. Pat. No. 5,282,624 discloses a golf club head having a face plate composed of a forged stainless steel material and having a thickness of 3 mm. Anderson, U.S. Pat. No. 5,344,140, for a Golf Club Head And Method Of Forming Same, also discloses use of a forged material for the face plate. The face plate of Anderson may be composed of several forged materials including steel, copper and titanium. The forged plate has a uniform thickness of between 0.090 and 0.130 inches.\nAnother invention directed toward forged materials in a club head is Su et al., U.S. Pat. No. 5,776,011 for a Golf Club Head. Su discloses a club head composed of three pieces with each piece composed of a forged material. The main objective of Su is to produce a club head with greater loft angle accuracy and reduce structural weaknesses. Aizawa, U.S. Pat. No. 5,346,216 for a Golf Club Head, discloses a face plate having a curved ball hitting surface.\nU.S. Pat. No. 6,146,571 to Vincent, et. al., discloses a method of manufacturing a golf club head wherein the walls are obtained by injecting a material such as plastic over an insert affixed to a meltable core. The core has a melt point lower than that of the injectable plastic material so that once the core is removed, an inner volume is maintained to form the inner cavity. The insert may comprise a resistance element for reinforcing the internal portion of the front wall of the shell upon removal of the core where the reinforcement element is comprised of aluminum with a laterally extending portion comprised of steel.\nU.S. Pat. No. 6,149,534 to Peters, et al., discloses a golf club head having upper and lower metal engagement surfaces formed along a single plane interface wherein the metal of the lower surface is heavier and more dense than the metal of the upper surface.\nU.S. Pat. Nos. 5,570,886 and 5,547,427 to Rigal, et al., disclose a golf club head of molded thermoplastic having a striking face defined by an impact-resistant metallic sealing element. The sealing element defines a front wall of the striking surface of the club head and extends upward and along the side of the impact surface to form a neck for attachment of the shaft to the club head. The sealing element preferably being between 2.5 and 5 mm in thickness.\nU.S. Pat. No. 5,425,538 to Vincent, et al., discloses a hollow golf club head having a steel shell and a composite striking surface composed of a number of stacked woven webs of fiber.\nU.S. Pat. No. 5,377,986 to Viollaz, et al., discloses a golf club head having a body composed of a series of metal plates and a hitting plate comprised of plastic or composite material wherein the hitting plate is imparted with a forwardly convex shape. Additionally, U.S. Pat. No. 5,310,185 to Viollaz, et al., discloses a hollow golf club head having a body composed of a series of metal plates, a metal support plate being located on the front hitting surface to which a hitting plate comprised of plastic or composite is attached. The metal support plate has a forwardly convex front plate associated with a forwardly convex rear plate of the hitting plate thereby forming a forwardly convex hitting surface.\nU.S. Pat. No. 5,106,094 to Desboilles, et al., discloses a golf club head having a metal striking face plate wherein the striking face plate is a separate unit attached to the golf club head with a quantity of filler material in the interior portion of the club head.\nU.S. Pat. No. 4,568,088 to Kurahashi discloses a wooden golf club head body reinforced by a mixture of wood-plastic composite material. The wood-plastic composite material being unevenly distributed such that a higher density in the range of between 5 and 15 mm lies adjacent to and extends substantially parallel with the front face of the club head.\nU.S. Pat. No. 4,021,047 to Mader discloses a golf club wherein the sole plate, face plate, heel, toe and hosel portions are formed as a unitary cast metal piece and wherein a wood or composite crown is attached to this unitary piece thereby forming a hollow chamber in the club head.\nU.S. Pat. No. 5,624,331 to Lo, et al. discloses a hollow metal golf club head where the metal casing of the head is composed of at least two openings. The head also contains a composite material disposed within the head where a portion of the composite material is located in the openings of the golf club head casing.\nU.S. Pat. No. 1,167,387 to Daniel discloses a hollow golf club head wherein the shell body is comprised of metal such as aluminum alloy and the face plate is comprised of a hard wood such as beech, persimmon or the like. The face plate is aligned such that the wood grain presents endwise at the striking plate.\nU.S. Pat. No. 3,692,306 to Glover discloses a golf club head having a bracket with sole and striking plates formed integrally thereon. At least one of the plates has an embedded elongate tube for securing a removably adjustable weight means.\nU.S. Pat. No. 5,410,798 to Lo discloses a method of manufacturing a composite golf club head using a metal casing to which a laminated member is inserted. A sheet of composite material is subsequently layered over the openings of the laminated member and metal casing to close off the openings in the top of both. An expansible pocket is then inserted into the hollow laminated member comprising sodium nitrite, ammonium chloride and water causing the member to attach integrally to the metal casing when the head is placed into a mold and heated.\nU.S. Pat. No. 4,877,249 to Thompson discloses a wood golf club head embodying a laminated upper surface and metallic sole surface having a keel. In order to reinforce the laminations and to keep the body from delaminating upon impact with an unusually hard object, a bolt is inserted through the crown of the club head where it is connected to the sole plate at the keel and tightened to compress the laminations.\nU.S. Pat. No. 3,897,066 to Belmont discloses a wooden golf club head having removably inserted weight adjustment members. The members are parallel to a central vertical axis running from the face section to the rear section of the club head and perpendicular to the crown to toe axis. The weight adjustment members may be held in place by the use of capsules filled with polyurethane resin, which can also be used to form the faceplate. The capsules have openings on a rear surface of the club head with covers to provide access to adjust the weight means.\nU.S. Pat. No. 2,750,194 to Clark discloses a wooden golf club head with weight adjustment means. The golf club head includes a tray member with sides and bottom for holding the weight adjustment preferably cast or formed integrally with the heel plate. The heel plate with attached weight member is inserted into the head of the golf club via an opening.\nU.S. Pat. No. 5,193,811 to Okumoto, et al. discloses a wood type club head body comprised primarily of a synthetic resin and a metallic sole plate. The metallic sole plate has on its surface for bonding with the head body integrally formed members comprising a hosel on the heel side, weights on the toe and rear sides and a beam connecting the weights and hosel. Additionally, U.S. Pat. No. 5,516,107 to Okumoto, et al., discloses a golf club head having an outer shell, preferably comprised of synthetic resin, and metal weight member/s located on the interior of the club head. A foamable material is injected into the hollow interior of the club to form the core. Once the foamable material has been injected and the sole plate is attached, the club head is heated to cause the foamable material to expand thus holding the weight member/s in position in recess/es located in toe, heel and/or back side regions by pushing the weight member into the inner surface of the outer shell.\nU.S. Pat. No. 4,872,685 to Sun discloses a wood type golf club head wherein a female unit is mated with a male unit to form a unitary golf club head. The female unit comprises the upper portion of the golf club head and is preferably composed of plastic, alloy, or wood. The male unit includes the structural portions of sole plate, a face insert consists of the striking plate and weighting elements. The male unit has a substantially greater weight being preferably composed of a light metal alloy. The units are mated or held together by bonding and or mechanical means.\nU.S. Pat. No. 5,398,935 to Katayama discloses a wood golf club head having a striking face wherein the height of the striking face at a toe end of the golf club head is nearly equal to or greater than the height of the striking face at the center of the club head.\nU.S. Pat. No. 1,780,625 to Mattern discloses a club head with a rear portion composed of a light-weight metal such as magnesium. U.S. Pat. No. 1,638,916 to Butchart discloses a golf club with a balancing member composed of persimmon or a similar wood material, and a shell-like body composed of aluminum attached to the balancing member.\nThe Rules of Golf, established and interpreted by the United States Golf Association (“USGA”) and The Royal and Ancient Golf Club of Saint Andrews, set forth certain requirements for a golf club head. The requirements for a golf club head are found in Rule 4 and Appendix II. A complete description of the Rules of Golf are available on the USGA web page at www.usga.org. Although the Rules of Golf do not expressly state specific parameters for a golf club face, Rule 4-1e prohibits the face from having the effect at impact of a spring with a golf ball. In 1998, the USGA adopted a test procedure pursuant to Rule 4-1e which measures club face COR. This USGA test procedure, as well as procedures like it, may be used to measure club face COR.\nAlthough the prior art has disclosed many variations of multiple material club heads, the prior art has failed to provide a multiple material club head with a high coefficient of restitution and greater forgiveness for the typical golfer."}
{"text": "The invention relates to a clutch disc for a friction clutch, in particular of a motor vehicle.\nFrom DE-A-3,340,896(U.S. Pat. No. 4,635,780) a friction clutch disc is known which comprises a hub with a hub disc which projects from the hub radially to an axis of rotation of the hub and with a cylindrical bearing face axially to the side of the hub disc and two lateral discs which are arranged axially on either side of the hub disc. The two lateral discs are rigidly connected to one another to form a unit together with the friction linings of the clutch disc and are rotatable relative to the hub disc about a limited angle of rotation. Each lateral disc has a central opening through which the hub passes. Several springs couple the lateral discs rotatably to the hub disc. A plastics mounting ring which is held non-rotatably in the central opening of a first lateral disc of the two lateral discs guides the two lateral discs on the hub. The mounting ring has the form of an annular disc which rests axially laterally on the lateral disc and from whose internal periphery spring tongues shaped to a cone extend through the central opening in the lateral disc. The spring tongues carry the lateral disc in their base region and lie with their ends under radial tension on a cylindrical bearing face of the clutch disc hub. A mounting ring of this type which can be produced inexpensively holds the lateral discs without clearance on the hub and compensates production tolerances.\nIt has been found that a mounting ring of the known type can compensate for errors of alignment between gearbox shaft and crankshaft only to a limited extent. In particular, with engine/gearbox constructions in which the gearbox shaft is no longer mounted in a pilot bearing of the crankshaft, as previously the case, relatively great errors of alignment between gearbox shaft and crankshaft can occur and can impair the operation of a torsional vibration damper provided in the clutch disc. Although the basic friction of the torsional vibration damper caused by the known mounting ring is relatively slight, it is still too great for a plurality of applications and detrimentally affects the damping devices of the torsional vibration damper, particularly if errors of alignment between the gearbox shaft and the crankshaft or inclination errors of these shafts occur. With the known clutch disc, the spring tongues lie with their tongue edges on the hub, leading to increased edge pressure, increased wear and differing coefficients of friction. A primary object of the invention is to provide a clutch disc provided with a torsional vibration damper, in which even relatively great errors of alignment between gearbox shaft and crankshaft can be compensated without impairing operation of the torsional vibration damper."}
{"text": "The present invention relates to a knife feed system in a rotary drum cutter provided in a tobacco shredder or the like.\nGenerally, in a shredding machine for tobacco, tobacco leaves as raw material are conveyed under compression to a shredding port by means of two upper and lower press conveyors, the shredding port being formed in front of the press conveyors, and are shredded with a rotary drum cutter rotating in close proximity to the shredding port.\nThe rotary drum cutter is composed of a rotary drum and a plurality of knives disposed at predetermined intervals on the outer peripheral surface of the rotary drum. The cutting edge of each knife is ground continually during operation of the cutter by means of a grinder provided separately so as to be suited to a high speed processing over a long period of time. Further for maintaining constant the spacing between the knife edge and the shredding port, a knife feed system for feeding each knife continuously or intermittently little by little is provided in the interior of the rotary drum.\nThe abrasion loss of the knife differs depending on the raw materials to be shredded. If the feed of the knife is small despite of a large abrasion loss, it will be impossible to attain sharpness of the knife even after grinding and so consequently difficult to obtain shredded tobacco leaves of good quality, causing an increased loss of raw material. Furthermore if the knife is ejected to a larger extent than necessary despite a small abrasion loss, the knife which is expensive will be wasted. Therefore, it is necessary that the feed of the knife be changed according to the raw material to be shredded. In this connection, the knife feed system in conventional rotary drum cutters is constructed so that each knife is ejected in interlock with the rotation of the rotary drum. This construction was given rise to the problem that the replacement of parts must be done for adjusting the amount of ejection of the knife according to the kind of raw material used. An additional problem has existed in that when changing the knife the spacing between the knife and the shredding port must be adjusted under no-load rotation of the drum without feed of raw material, thus requiring much time for the adjustment."}
{"text": "As compared with other permanent magnetic material, sintered NdFeB permanent magnetic material has some outstanding advantages such as high magnetic property and low cost, so that it has been widely developed and applied. At present, sintered NdFeB permanent magnetic material has been applied in many fields due to a relative higher comprehensive magnetic property.\nHowever, new energy and environmental protection have obtained increased concern and become a developing trend, and a permanent magnetic material with a high coercivity and a high remanence is required. A permanent magnetic material with a high coercivity needs relative more expensive elements dysprosium and/or terbium. But if too much such elements are added, neither the requirement of high remanence can be meet, nor the light weight of motor and high availability of electric energy and wind energy can be obtained.\nThe coercivity of prepared permanent magnet material at present has a significant difference from a theoretical limit 80 kOe, and a relative high content of Dy and/or Tb is needed in order to improve the coercivity of the permanent magnet material at present. In addition, when the coercivity of the permanent magnet material is increased, the remanence may be decreased. Therefore, it is required to improve the coercivity with only a small decrease of the remanence of the permanent magnetic material by using a small amount of Dy and/or Tb."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus and method for granulating plastic strands that emerge from nozzles in molten form and are supplied to a following, downward inclined runoff channel in which the plastic strands are transported by a liquid stream, in particular water stream, the apparatus further having a dewatering zone for the plastic strands and a granulator disposed after said dewatering zone for granulating the solidified plastic strands with the aid of at least one knife assembly.\n2. Description of the Related Art\nU.S. Pat. No. 5,182,115 discloses an apparatus for cooling, drying and granulating strands emerging from nozzles in molten form having a runoff channel disposed with its receiving-side end below the nozzles, and further having a device for producing a coolant stream on the runoff channel, while the runoff channel is followed by a granulator and furthermore a dewatering zone is disposed in the runoff channel before the discharge-side end, the runoff channel being provided in said dewatering zone with openings for free passage of the coolant.\nIn addition the runoff channel in this known apparatus is provided in its bottom in the area after the dewatering zone with admission nozzles for a gas or air stream over such a length and so close together that the strands are guided for removal of residual water along the runoff channel in largely frictionless fashion with respect to the bottom into the granulator disposed at its discharge-side end with a degree of dryness permitting immediate processing. This granulator has essentially two draw-in rolls for the plastic strands one above the other within a corresponding housing, further a knife roll cooperating with a corresponding counterknife. This knife roll and the counterknife granulate the plastic strands supplied by the draw-in rolls, and the granules thus produced by the granulator fall out through a discharge shaft for further processing.\nIn addition U.S. Pat. No. 4,180,539 discloses an apparatus of the above-defined type wherein the runoff channel is followed by an upward inclined dewatering and drying zone for the plastic strands that consists essentially of a circulating, endless and latticed conveyer belt on which the plastic strands leaving the runoff channel and thereby entraining at least part of the liquid stream are transported further toward the granulator. Above this circulating latticed conveyer belt a number of air nozzles are disposed for directing relatively sharp air jets (\"air knives\") onto the surface of the plastic strands. These air knives serve to dry the strands on their surfaces during transportation, whereby they blow away the water sticking thereto and remove it downward through the lattice belt. The latticed formation of the transport belt is necessary in the known apparatus of U.S. Pat. No. 4,180,539 in particular because relatively sharp air jets are directed onto the strands from above so that the air must likewise be removed downward through the lattice belt since the strands would otherwise not have the right contact pressure on the surface of the lattice belt; the strands subjected to the air knives would instead \"fly up\" and \"swirl around\" as it were. In addition, the known dewatering and drying zone ensures only a dewatering and drying process with respect to the strand surfaces while the interior of the strands still has an elevated temperature compared to the surface, so that the plastic strands entering the granulator directly after the dewatering and drying zone are not yet equalized in their temperature profile.\nIt is fundamentally necessary, however, that the plastic strands to be processed into granules be neither too soft nor too brittle when they are supplied to the knife roll of the granulator. In view of this requirement and also because of the extremely different properties of the plastics marketed in granular form, these properties ranging from \"very elastic\" to \"extremely brittle\", the cooling devices used between an extruder and a granulator in known apparatuses are already adapted by the manufacturer to the type of plastic later to be processed in order to ensure that the newly extruded plastic strands are cooled with such an intensity that they have the desired strength properties when they reach the knife roll of the granulator. When processing certain plastics, in particular glass fiber-reinforced plastics, one requires that the granulating temperature be over about 100.degree. C., whereby the temperature profile should be as homogeneous as possible across the strand cross section. The same requirement also holds quite generally for example in the processing of plastics that are fundamentally relatively poor heat conductors.\nIn the apparatus known from U.S. Pat. No. 4,180,539 the dewatering and drying process does not permit the plastic strands to be given a homogeneous temperature distribution across the entire strand cross section before they enter the granulator, because the dewatering and drying process concentrates in particular on the strand surfaces."}
{"text": "Regular expression search operations are employed in various applications including, for example, intrusion detection systems (IDS), virus protections, policy-based routing functions, internet and text search operations, document comparisons, and so on. A regular expression can simply be a word, a phrase or a string of characters. For example, a regular expression including the string “gauss” would match data containing gauss, gaussian, degauss, etc. More complex regular expressions include metacharacters that provide certain rules for performing the match. Some common metacharacters are the wildcard “.”, the alternation symbol “|”, and the character class symbol “[ ]”. Regular expressions can also include quantifiers such as “*” to match 0 or more times, “+” to match 1 or more times, “?” to match 0 or 1 times, {n} to match exactly n times, {n,} to match at least n times, and {n,m} to match at least n times but no more than m times. For example, the regular expression “a.{2}b” will match any input string that includes the character “a” followed exactly 2 instances of any character followed by the character “b” including, for example, the input strings “abbb,” “adgb,” “a7yb,” “aaab,” and so on.\nTraditionally, regular expression searches have been performed using software programs executed by one or more processors, for example, associated with a network search engine. For example, one conventional search technique that can be used to search an input string of characters for multiple patterns is the Aho-Corasick (AC) algorithm. The AC algorithm locates all occurrences of a number of patterns in the input string by constructing a finite state machine that embodies the patterns. More specifically, the AC algorithm constructs the finite state machine in three pre-processing stages commonly referred to as the goto stage, the failure stage, and the next stage. In the goto stage, a deterministic finite state automaton (DFA) or search tree is constructed for a given set of patterns. The DFA constructed in the goto stage includes various states for an input string, and transitions between the states based on characters of the input string. Each transition between states in the DFA is based on a single character of the input string. The failure and next stages add additional transitions between the states of the DFA to ensure that a string of length n can be searched in exactly n cycles. More specifically, the failure and next transitions allow the state machine to transition from one branch of the tree to another branch that is the next best (i.e., the longest prefix) match in the DFA. Once the pre-processing stages have been performed, the DFA can then be used to search any target for all of the patterns in the pattern set.\nOne problem with prior string search engines utilizing the AC algorithm is that they are not well suited for performing wildcard or inexact pattern matching. As a result, some search engines complement the AC search technique with a non-deterministic finite automaton (NFA) engine that is better suited to search input strings for inexact patterns, particularly those that include quantifiers such as “k” to match 0 or more times, “+” to match 1 or more times, “?” to match 0 or 1 times, {n} to match exactly n times, {n,} to match at least n times, and {n,m} to match at least n times but no more than m times.\nEmploying an NFA engine to search an input string for a regular expression generally involves converting the regular expression into an NFA search tree that includes a number of states interconnected by goto or “success” transitions. Then, to search the input string for the regular expression, a state machine starts at an initial state of the NFA and transitions to one or more states of the NFA according to its goto transitions. If in a given state the input character matches the goto transition, the goto transition is taken to the next state in the string path and a cursor is incremented to point to the next character in the input string. Otherwise, if there is not a character match at a particular state, the state becomes inactive.\nFor example, FIG. 1 shows an NFA 100 that embodies the regular expression R1=“[a-z][a-z0-9]{5},” which includes a prefix pattern P1 including the character class “[a-z]” followed by a quantified character class “[a-z0-9]{5}.” An input string of characters will match the regular expression R1 if the input string includes a first character that matches the prefix character class [a-z] followed by n={5} characters that match the quantified character class [a-z0-9]. The NFA 100 for R1 includes an initial state or root node S0, intermediate states S1-S5, and a match state. The sequence of states S0-S5-Match are connected by goto transitions representing character matches with an input string, as indicated in FIG. 1. For example, if an input character within the prefix character class [a-z] is received while the state machine is in the initial state S0, the state machine transitions from S0 to S1 along the “[a-z]” goto transition, and the cursor is incremented to the next input character in the input string. Then, if the next input character is within the quantified character class [a-z0-9], the state machine transitions from S1 to S2 along the “[a-z0-9]” goto transition. Then, if the next input character is within the quantified character class [a-z0-9], the state machine transitions from S2 to S3 along the “[a-z0-9]” goto transition, and so on until a match is detected at the match state. Conversely, if at S1-S5 the input character is anything other than a character within the quantified character class [a-z0-9], then no transition occurs and the current state is deactivated. The initial state S0 remains active for the NFA, as indicated by the arrow 101, which matches on the wildcard “.”. Thus, regardless of what the input character is, the root state S0 remains active.\nThe regular expression R1=“[a-z][a-z0-9]{5}” is considered a complex regular expression because multiple states of the corresponding NFA 100 can be active at the same time. More specifically, because the quantified character class [a-z0-9] overlaps with (i.e., is a superset of) the prefix character class [a-z], each input character that matches the quantified character class also matches the prefix character class, and therefore not only causes the state machine to transition from one of states S1 to S5 to another of states S2-Match but also causes the state machine to activate an additional instance of state S1. The activation of the additional instance of S1 indicates the beginning of another separate and overlapping portion of the input string that can potentially match the regular expression R1. For example, Table 1 depicts a search operation between an input string IN1=“abcdefgh” and the regular expression R1 according to the NFA 100 of FIG. 1.\nTABLE 1CycleInput CharacterActive StatesAction1a02b1,03c2,1,04d3,2,1,05e4,3,2,1,06f5,4,3,2,1,0match7g5,4,3,2,1,0match8h5,4,3,2,1,0match\nAs shown in Table 1, the number of active states in the NFA 100 can quickly escalate because after the match between the first input character “a” and the prefix character class [a-z], each subsequent input character that matches both the prefix character class [a-z] and the quantified character class [a-z0-9] triggers another potentially matching sub-string and counts towards the quantified number {5} of matches of the character class [a-z0-9], which in turn can quickly exhaust counter resources within the NFA engine.\nUnfortunately, the available memory to support an NFA engine is limited, and therefore the number of active states that can be maintained by the NFA engine is limited. As a result, information regarding active states may become lost due to unavailable state memory. Moreover, increasing the number of active states for each regular expression search operation may undesirably degrade search performance by overly taxing processing resources during each compare cycle.\nThus, a need exists for an NFA-based search system that minimizes the memory resources devoted to maintaining lists of active states and tracking numerous overlapping sub-string matches.\nLike reference numerals refer to corresponding parts throughout the drawing figures."}
{"text": "1. Field of the Invention\nThis invention relates to the art of label roll packages.\n2. Brief Description of the Prior Art\nThe following U.S. patents are made of record:\n______________________________________ 354,984 J. Keleher Dec. 28, 1886 2,748,931 J. Taylor June 5, 1956 2,912,102 A. Scott Nov. 10, 1959 3,161,999 G. Klusmire Dec. 22, 1964 3,338,399 R. Burt Aug. 29, 1967 3,562,999 R. Barbedienne Feb. 16, 1971 ______________________________________\nThe Keleher patent discloses a container, a compound braid roll and a case therefor. The Taylor patent discloses a dispensing container for a plurality of rolls. The Scott patent discloses a package for spools wrapped with a transparent plastics material. The Klusmire patent discloses a tray of objects overwrapped with plastics film material in which a shaping tool is used for forcing the material to at least partially conform to the contents. The Burt patent discloses a package wrapped with a heat shrinkable film which is subjected to heat to cause the overwrap film to shrink snugly around the frame to increase the rigidity of the frame structure. The Barbedienne patent discloses a container for packing flexible tubes in an envelope that is contracted by heat."}
{"text": "1. Field of Invention\nThe invention relates generally to the field of oil and gas production. More specifically, the present invention relates to a perforating system. Yet more specifically, the present invention relates to a shaped charge having a modified boss for use in a perforating gun system.\n2. Description of Prior Art\nPerforating systems are used for the purpose, among others, of making hydraulic communication passages, called perforations, in wellbores drilled through earth formations so that predetermined zones of the earth formations can be hydraulically connected to the wellbore. Perforations are needed because wellbores are typically completed by coaxially inserting a pipe or casing into the wellbore. The casing is retained in the wellbore by pumping cement into the annular space between the wellbore and the casing. The cemented casing is provided in the wellbore for the specific purpose of hydraulically isolating from each other the various earth formations penetrated by the wellbore.\nPerforating systems typically comprise one or more perforating guns strung together, these strings of guns can sometimes surpass a thousand feet of perforating length. In FIG. 1 an example of a perforating system 4 is shown. For the sake of clarity, the system 4 depicted comprises a single perforating gun 6 instead of a multitude of guns. The gun 6 is shown disposed within a wellbore 1 on a wireline 5. The perforating system 4 as shown also includes a service truck 7 on the surface 9, where in addition to providing a raising and lowering means, the wireline 5 also provides communication and control connectivity between the truck 7 and the perforating gun 6. The wireline 5 is threaded through pulleys 3 supported above the wellbore 1.\nPerforating guns typically include a cylindrical gun strip 16 coaxially housed within a gun body 14. Shaped charges 8 are provided within the gun strip 16 and aimed generally perpendicular to the axis of the wellbore 1. FIG. 2 provides an example of a shaped charge 8 that includes a housing 18, a liner 22, a quantity of high explosive 24 inserted between the liner 22 and the housing 18, and a booster charge 26 adjacent the base of the high explosive 24.\nThe shaped charges 8 are typically connected to a detonating cord 27 which is affixed to the shaped charge 8 by a case slot 25 proximate to the booster charge 26. Igniting the detonation cord 27 creates a compressive pressure wave along its length that initiates the booster charge 26 that in turn ignites the high explosive 24. When the high explosive 24 is detonated, the force of the detonation collapses the liner 22 and ejects it from one end of the charge 8 at very high velocity in a pattern called a “jet” 12. The jet 12 perforates the casing and cement lining the wellbore 1 and creates a perforation 10 that extends into the surrounding formation 2.\nShaped charges 8 also include a boss 20 protruding outward from the case 18 perpendicular to the axis ASC of the case 18. The boss 20 fully circumscribes the case 18 outer circumference. A perspective example of the gun strip 16 is provided in FIG. 3 illustrating holes 28 formed through the gun strip 16 for receiving shaped charges 8 therein. The shaped charge 8 is inserted into the hole 28 until the boss 20, whose outer diameter extends past the hole 28 outer diameter, contacts the outer surface of the gun strip 16. Thus the boss 20 supports the shaped charge 8 in the gun strip 16 and vertically aligns the shaped charge 8 in the gun strip 16. However, because the boss 20 is generally planar but the gun strip 16 outer diameter is curvilinear, the entire radius of the boss 20 does not lie above the hole 28, but instead the hole 28 outer diameter is shaped to accommodate the shaped charge 8 placement. Accordingly although the shaped charge 8 outer diameter is substantially circular, the typical gun tube 16 hole 28 is not. This can be a problem in certain perforating guns employing a “high density” shot pattern. For example, FIG. 3 illustrates an example of a gun tube 16 having high density shot pattern wherein adjacent holes 28 are sufficiently close that a web portion 30 between the holes 28 is too narrow to provide adequate structural support."}
{"text": "Feature detection of a two-dimensional code is a first critical step in two-dimensional data processing. This is because it is impossible to continue the processing once the detection is failed. However, a two-dimensional code generally has a large amount of data, hence whether the two-dimensional code feature can be quickly detected has a direct influence on the utilization performance of the two-dimensional code. Therefore, a quick and accurate detection of two-dimensional code feature has become an important demand in the two-dimensional code data processing.\nA quick response (QR) code is a type of two-dimensional code. The QR code has three position detection patterns which are located at an upper-left corner, an upper-right corner and a lower-left corner of the QR code respectively. Each of the position detection patterns may be considered to be formed by three overlapped concentric squares which are a 7*7 dark color module, a 5*5 light color module and a 3*3 sized dark color module, respectively. Therefore, a proportion of lengths of black and white line segments divided by the center of the position detection pattern meets 1:1:3:1:1.\nAt present, in most two-dimensional code feature detection methods, scan and detection are performed row by row, and the detection speed is slow.\nChinese invention patent No. CN104517090A, discloses a method and system for detecting a QR code detection pattern. The method includes: detecting a QR code in a first direction, detecting a first line segment of which a proportion of lengths of black: white: black: white: black line segments meets a preset proportion so as to determine a center of the first line segment; detecting, in a second direction, a straight line passing a center point of the first line segment, to obtain a second line segment of which a proportion of lengths of black: white: black: white: black line segments meets a preset proportion, and determining a center point of the second line segment as a center point of the detection pattern. Although only a row-by-row full coverage detection is required in the first direction detection and the detection speed is improved, only the center point of the detection pattern can be found in detecting the detection pattern of a QR code image, and the boundary lines of the detection pattern of the complete QR code can not be determined. Particularly, the boundary lines of the detection pattern of the QR code can not be determined accurately in the case that the QR code is deformed and distorted."}
{"text": "1. Field of the Invention\nThe present invention relates to digital information storage, and more particularly, to a method for adjusting a delay time for opening a servo sector detection window in a disk drive.\n2. Description of the Prior Art\nA magnetic disk drive provides a nonvolatile disk media for storage of user and application information. The disk media is generally formatted to store the information in concentric rings or tracks. The concentric tracks are followed using embedded servo wedges that are sampled during disk access operations. To switch from one track to another, the disk drive must perform a mechanical seek operation to move a transducer head to the desired track. During a seek, the stored in formation may not be accessed. When the disk drive is providing streaming media information, such as a video or audio information, to a host, the information must be buffered to prevent annoying discontinuities during playback of the streaming media. Also, to increase the storage capacity of the disk drive, the disk media may be divided into zones, with more uniformly-sized data blocks being written between the servo wedges on the outer tracks than are written between the servo wedges on the inner tracks. Nevertheless, the linear circumferential distance of a track wires as its radial distance from the center of the disk varies. Thus, a substantial number of the tracks often have a portion that is unused because it is smaller that the uniformly-sized data blocks. Alternatively, the tracks may have data blocks that straddle servo wedges using techniques requiring increased format overhead.\nAccordingly, there exists a need for a disk drive and related techniques for using a magnetic disk media format that efficiently stores information and that allows more continuous information playback from a magnetic disk. The present invention satisfies this need."}
{"text": "1. Field of the Invention\nThis invention relates generally to ink compositions and, more particularly, this invention relates to ink compositions useful in jet printing applications with plain paper at room temperature, and methods of ink jet printing on plain paper.\n2. Description of Related Art\nPrinting according to the \"ink jet printing\" principle, and apparatus for carrying out such printing operations, are well known. In general terms, a fluid ink is forced, under pressure and often at an elevated temperature, through a very small orifice in a printing head.\nIn so-called \"continuous\" jet printing operations, ink droplets are passed through a charging area wherein individual droplets receive an electrical charge in response to a signal. The droplets then pass through an electrical field, causing a varied deflection of the individual droplets dependent on the intensity of the charge and field. The droplets are produced continuously, and are directed to the substrate to be printed or, alternatively, to a by-pass gutter.\nDue to the nature of the \"continuous\" jet printing process, inks used therein must be conductive so as to accept a charge. Such inks are generally water-based, as aqueous solutions or dispersions are easily rendered conductive. Water-based inks are generally characterized as having a low viscosity and high volatility. In addition, the interaction of water with paper can cause dimensional changes in the paper. These characteristics can be disadvantageous.\nIf the viscosity of a jet printing ink is excessively low, misting and the production of stray droplets can result, which in turn result in poor print quality. High volatility of a jet printing ink can result in clogging of jet nozzles due to evaporation. Also, only a few water-based inks are capable of use on plain paper (i.e. a paper not bearing a pigment coating) and are capable of providing a high dye concentration on the surface of the paper so that it is possible to produce fine-sized droplets and to avoid build-up of dyestuff on the substrate surface while providing high color intensity.\nSo-called \"drop on demand\" (DOD) systems differ from continuous jet printing systems in that ink droplets are expelled from the nozzle of a printing head only when required during the printing process.\nIn a DOD printing system, ink is provided from a reservoir via a supply system to a nozzle/actuator system in the printing head. The actuators are generally of the piezoelectric or bubble type, and pressurize ink upon activation to force the ink from a chamber defined by the actuator through a jet nozzle associated therewith.\nSince inks used in DOD systems need not be conductive, they have generally been formulated using primarily non-aqueous solvents such as ethylene glycols, particularly diethylene glycols, which are characterized as having low volatility. As such, ethylene glycol and similar solvent-based inks do not evaporate quickly and, thus, exhibit reduced clogging problems.\nHowever, ethylene glycol and other non-aqueous solvents have a relatively high viscosity (which varies significantly with temperature) and thus require relatively more energy for production of droplets.\nSolvent-based inks exhibit desirable drying characteristics and are generally biologically stable while water-based inks generally tend to sorb air from the environment and also are susceptible to biological fouling.\nIn some cases, it may be desirable to use a non-conductive ink (or one having extremely low conductivity) in DOD printers. For example, in DOD actuators wherein piezoelectric electrodes are disposed within an ink-containing chamber, the ink is preferably non-conductive in order to minimize or eliminate electrochemical effects.\nAnother class of inks used in jet printing are referred to in the art as \"hot melt\" inks and comprise wax or other thermoplastic materials such that the ink is solid or semisolid at ambient temperatures but fluid at elevated temperatures. The heated ink fluid solidifies when it comes in contact with a paper target.\nSuch inks generally tend to exhibit excellent dye stability, lightfastness, and no solvent volatility. However, these inks accommodate only relatively low dye concentration and therefore necessitate the production of relatively large drops in order to provide acceptable color intensity. Also, the thermoplastic component of such inks tends to form raised waxy deposits which are easily smeared or scratched off by rubbing.\nHot melt jet inks tend to be viscous and thus require high energy input in order to eject the drop from the printer nozzle. Also, since the entire ink must be rendered fluid by heating, high energy inputs and lengthy time periods are required to reach the printer operating temperature."}
{"text": "The present invention relates generally to the field of computer applications, and more particularly to improved performance of mobile device applications.\nA mobile application is a computer program designed to run on mobile devices such as smartphones and tablet computers. Mobile application development for mobile devices requires consideration of the constraints of mobile devices as compared with personal computers. Mobile devices typically have less powerful processors than personal computers and have user interface constraints such as a smaller screen. For access and interaction with enterprise systems, mobile device user interfaces rely on back-ends to support data routing, security, authentication, and authorization. This functionality is typically supported by various middleware components such as application servers."}
{"text": "The present invention relates to a method and device suitable for storing and transporting animal semen, and more particularly to a method and disposable device for storing and transporting equine semen.\nA device for transporting equine and other animal semen was introduced to the equine breeding market in 1984 with the trade name xe2x80x9cThe Equitainerxe2x80x9d (at least a partial description of which can be found in U.S. Pat. No. 4,530,816). Since that time, horse breeding by artificial insemination using transported chilled semen has become dominant in the equine breeding sector. The technique of slowly cooling the sample en route to its destination has been an important contribution.\nThe Equitainer is carefully engineered to give a repeatable and reliable cooling rate and final temperature that is substantially independent of external conditions. Although the performance of the Equitainer is good, the structure is relatively bulky and not inexpensive to produce. The cost of the Equitainer implies that the container must be returned to the shipper after use. In many cases, this can be inconvenient.\nThe conventional Equitainer achieves a required cooling rate by placing a sample in a conductive cavity, known as the Isothermalizer, and interposing a barrier between the sample and the coolant. This configuration determines the rate at which heat can be transferred between the coolant and the sample. The volume of the sample is then adjusted using thermal ballast, to result in a desired controlled cooling rate. The sample and coolant assembly is packaged inside a large insulated cylindrical container, and dispatched to the user. There is an initial cooling rate when the start temperature is at about 37xc2x0 C. of close to xe2x88x920.3xc2x0 C./minute. The cooling continues slowly en route, down to a minimum temperature typically 6xc2x0 C., depending on exterior conditions.\nIt has been shown and reported in the industry that the best results are obtained for horse semen fertility when the cooling rate at 20xc2x0 C. is in the vicinity of xe2x88x920.1xc2x0 C./min. This is to avoid a high rate of change to the cell temperature at the temperature where the lipid cell membranes are changing phase. The low temperature rate of change likely allows the proteins in the cell surface time to crystallize in an orderly fashion and avoid cell damage. The Equitainer design was based on this finding.\nIn addition to the Equitainer, a number of other conventional devices use a process of convection to cool the sample contained within the device. With convection, cooled air transfers heat from the sample to the coolant. The sample is normally transported in syringes ready for insemination, and a thermal barrier is introduced between the coolant and the sample to adjust cooling rate. The thermal barrier used is typically a Styrofoam sheet with one or more notches cut out. The Styrofoam sheet is placed between coolant block and sample, so that air can flow through the notches as it convects. The coolant block is normally placed above the sample. Cool air then falls through the holes in the Styrofoam and cools the sample. The sample is placed in a foamed plastic box, next to the Styrofoam thermal barrier and the coolant block.\nThe convection system does not independently control cooling rate and final temperature. In addition, the convection system is sensitive to orientation. This is evidenced by the fact that the cooling rate is about twice as fast when the system is upright verses when it is inverted, representing the difference between air convection cooling and air conduction cooling. To make a system in which the cooling rate is independent of orientation, convective cooling cannot be used. The failure to maintain a uniformly controlled cooling rate can have negative effects on the success rates of equine insemination.\nThere is a need in the art for a disposable container with similar reliable cooling rate properties as the Equitainer, but with the added constraint that the materials and components required must not only produce accurate and repeatable results, but must be simple, disposable, and inexpensive. The present invention is directed toward further solutions to address this need.\nIn accordance with one embodiment of the present invention, a storage and transportation container for storing and transporting a semen sample is provided. The container includes an insulated box having at least a first portion being a box protruding ridge, at least a second portion being a box recessed groove, and a floor. An insulated cover is provided having at least a first portion having a cover recessed groove and being located to receive the box protruding ridge only when the insulated cover is placed on the insulated box in a predetermined orientation. The insulated cover also has at least a second portion having a cover protruding ridge and being located to be received into the box recessed groove only when the insulated cover is placed on the insulated box in the predetermined orientation. A cooling pack is disposed on the floor of the insulated box. An insert is disposed to rest on the cooling pack inside the insulated box, such that the cooling pack is beneath the insert when the insulated box is in an upright position. The insert has at least one chamber for receiving a storage capsule holding the semen sample. The insulated cover frictionally fits on the insulated box and makes contact with the insert to hold the insert in place. The insert makes sufficient contact with the cooling pack to enable conductive cooling of the semen sample in the insert.\nIn accordance with various aspects of the present invention, the insulated box and the insulated cover are formed of foam plastic, such as Styrofoam, or the like. Likewise, the insert can also be made of foamed plastic, such as polyethylene, or the like. The choice of material rests on insulation properties and relative cost.\nIn accordance with further aspects of the present invention, the at least one chamber is sized and dimensioned to hold a storage capsule suitable for holding the semen sample. The semen sample is equine semen. The storage capsule is at least one of a syringe and a tube.\nThe storage capsule can be sized and dimensioned differently, and can include sizes with volume capacities of 20 ml or 50 ml. The relative size of the storage capsule will affect the number of chambers required in the insert to hold the storage capsules. Since the total amount of material to be cooled will affect the cooling rate, thermal ballast can be added to the storage capsule, if necessary, to keep the total volume in the range of 80 ml to 120 ml in one example embodiment. The thermal ballast is typically water, contained in a suitable tube, initially at substantially the same temperature as the semen sample. Smaller capsules would result in more capsules being required and thus a plurality of chambers required in the insert. A sum of the capacity of the semen sample and the thermal ballast able to be stored in the at least one chamber according to one embodiment is about 90 ml to about 110 ml, or about 80 ml to about 120 ml. With such a semen sample quantity, a cooling rate of the semen sample in the insert is about 0.1xc2x0 C./min at about 20xc2x0 C.\nIn accordance with further aspects of the present invention a method of packing a semen sample for storage or transportation includes providing an insulated box having a floor. A cooling pack is placed on the floor of the insulated box. An insert is placed on top of and substantially in contact with the cooling pack, the insert having at least one storage capsule holding the semen sample disposed within at least one chamber of the insert. An insulated cover is placed on top of and in contact with the insert, the insulated cover having a friction fit with the insulated box."}
{"text": "Neurostimulation (NS) systems are devices that generate electrical pulses and deliver the pulses to nervous tissue to treat a variety of disorders. For example, spinal cord stimulation (SCS) has been used to treat chronic and intractable pain. Another example is deep brain stimulation (DBS), which has been used to treat movement disorders such as Parkinson's disease and affective disorders such as depression. SCS therapy, delivered via epidurally implanted electrodes, is a widely used treatment for chronic intractable neuropathic pain of different origins. Traditional tonic therapy evokes paresthesia covering painful areas of a patient. During SCS therapy calibration, the paresthesia is identified and localized to the painful areas by the patient in connection with determining correct electrode placement.\nRecently, new stimulation therapies such as burst stimulation and high frequency stimulation, have been developed, in which closely spaced high frequency pulses are delivered. In general, conventional neurostimulation systems seek to manage pain and other pathologic or physiologic disorders through stimulation of select nerve fibers that carry pain related signals. However, nerve fibers and brain tissue carry other types of signals, not simply pain related signals.\nAlthough some neurological disorders have been treated through known neurostimulation methods, many other neurological disorders exhibit physiological complexity, functional complexity, or other complexity and have not been adequately treated through known neurostimulation methods."}
{"text": "1. Field of the Invention:\nThis invention relates generally to lead frames which are used to constitute various electronic components such as IC's or LSI's, and more particulary to such a lead frame having a pin holding structure which serves to prevent undesirable deformation of the lead pins during transport of the lead frame or during assembly thereof into an electronic component. The invention also relates to a pin holding method for a lead frame.\n2. Description of the Prior Art:\nAs is well known in the art, lead frames are manufactured by photoetching or press-working a thin metallic strip into a predetermined form to have a longitudinally spaced islands and a plurality of lead pins associated with each island. The metallic strip is usually made of Fe-Ni based alloy, copper, or copper based alloy.\nDue to recent demand for further integration or sophistication of semiconductor devices, lead frames to be incorporated in such devices are required to have a very small thickness with its lead pins slenderized and densely arranged. Therefore, the lead pins are becoming increasingly liable to deformation during transport of the lead frame and/or during assembly thereof into an IC, particularly during chip bonding and wire bonding. Such deformation of the lead pins results in poor quality and low yield.\nIn order to eliminate the above problem, it has been proposed to apply a heat-resistant adhesive tape onto the inner lead pins of a lead frame, so that the inner lead pins are prevented from deformation and resultant shorting. More specifically, the tape consists of a polyimide film coated at one side with an acrylic adhesive. Such adhesive tape is available under the tradename \"PYRALUX\" from Du Pont Company.\nIn use, a blank adhesive tape which is originally wide is punched by a die into a segment of suitable shape and dimensions depending on the particular type of lead frame being processed. The tape segment is then transferred onto a specified surface portion of the lead frame, and the taped lead frame is subjected to heat and pressure to maximize adhesion of the tape.\nThe above pin holding tape, however, has a disadvantage that it necessitates the use of a costly punching die which is designed to suit a particular type (length, width, configuration, etc.) of lead frame and which must be replaced in a time-taking operation by a differetly designed die to suit another type of lead frame. After punching, the tape produces a substantial residual portion, resulting in economical problem in view of the expensive materials of which the tape is made. Further, the tape has a possibility of coming off relatively easily from the lead frame surface."}
{"text": "Aerated frozen products generally require the mixing of selected liquid ingredients with a prescribed volume of air and freezing of the resultant mixture and dispensing of the finished product. The desirability of the finished product is often related directly to the manner and to the degree in which the air is metered and blended with the liquid ingredients of the mixture and the manner in which the blended mix is frozen. The prior art is replete with examples of apparatus for dispensing ice cream and other semifrozen dairy products such as soft ice cream and frozen yogurt.\nUsually in such dispensers, the liquid mix containing, in the case of ice cream, ingredients such as cream, milk, condensed milk, syrup, etc. is delivered to a freezing chamber containing air or another non-toxic gas. In the freezing chamber, the mixture is slowly frozen and mechanically agitated by blades or the like to incorporate the air or gas in the mixture. This aeration is especially important and produces an increase in bulk corresponding to a decrease in the product specific gravity, thereby providing what is referred to in the ice cream industry as \"over-run\". The slow freezing and continuous agitation of results in the formation of ice crystals and particles of varying size which can detract from the palatability of the resultant product.\nThe chamber typically serves both as the means of freezing and as a reservoir from which individual servings are withdrawn. When the ice cream is partially frozen to the proper consistency, it is pumped or extruded from the freezing chamber to the dispenser outlet for delivery, on demand, to containers such as cups and cones. Examples of such apparatus are disclosed in U.S. Pat. Nos. 3,904,085; 3,954,126 and 4,201,588.\nThe prior apparatus do not necessarily achieve close control of over-run, which is a prime factor affecting product palatability and profit in the manufacture of ice cream and similar products. If there is not enough aeration in the finished product, not only may the product be so dense as to be unpalatable, but also more liquid mix is required to make a given volume of the product, thereby lowering profit. On the other hand, if there is too much air in the finished product, the product may be considered too \"fluffy\" and of lower quality and value. Also, in many states there are regulations against providing excessive over-run in certain products to protect the consumer from being charged for a product that consists largely of air.\nConventional dispensers are usually dedicated to dispensing one or two flavors of product and, in some cases, a combination (\"twist\") of the two. For example, in an ice cream shop, there may be one machine with two separate freezing chambers for making and dispensing chocolate and vanilla ice cream, a second two-chamber machine for making and dispensing strawberry and banana ice cream, a third machine dedicated to making and dispensing coffee and frozen pudding flavors, and so on. The reason for this is that each chamber typically contains a volume of ice cream greater than is required for a single serving. In order to dispense a different flavor ice cream, that chamber must be emptied before the new flavor can be made in that chamber and appear at the outlet of the dispenser; additionally, the vat of preflavored mix from which the material being aerated and frozen is drawn must also be cleaned. While high volume ice cream shops and confectionery stores may have sales to justify the presence of several dispensing machines dispensing many different products, smaller sales outlets can usually only afford one or two such machines and are thus restricted in the number of flavors that they can offer to customers.\nFurther, because the product is typically formed in a quantity that is greater than that to be dispensed at any one serving, the excess product remains in the chamber after formation and until the next serving is required. The excess is thus subjected to further mechanical beating by any blades in the chamber, as well as to excess freezing which promotes crystallization. Because of the quantity of the premixed flavors, and the continuous freezing and beating of several quarts of the product, the freshness and palatability of the product may be adversely affected in stores with slow sales of the product.\nAnother disadvantage of the prior dispensers is that they have many interior surfaces and moving parts that are difficult and time-consuming to clean and to maintain. At the end of each day or at other intervals prescribed by local Health Department regulations, each dispenser must be purged of any remaining product, and its chamber walls, pumps, and other internal parts cleaned thoroughly to prevent growth of bacteria that could contaminate product being delivered by the dispenser. Not only is the cleaning operation expensive in terms of downtime, it is also costly in terms of product waste and is an unpleasant and difficult job to get employees to do properly.\nThere have been attempts to produce and dispense aerated frozen products on a continuous basis by atomization. In one example of this type of dispenser, disclosed in U.S. Pat. No. 2,594,422, a liquid mix and a liquid refrigerant are mixed in a mixing chamber and fed through a rotary emulsifier which forms a single liquid emulsion. The emulsion is then jetted from a spray nozzle, into a separating chamber where the liquid refrigerant evaporates, effectively freezing adjacent droplets of mix and thereby transforming them into flakes or a fine powder. Groups of flakes then agglomerate into larger particles containing an interior void. The refrigerant adsorbed on the interior surfaces of the void may further expand in the separating chamber to increase the bulk of the particles and thus decrease their density. The liquid refrigerant evaporated in the separating chamber is drawn off and reused. The frozen particles fall into the chamber of a screw-type extruder which compresses the particles at a controlled rate to give the final product the desired density and drives the frozen product to the dispenser outlet.\nThe over-run is formed by, and limited by, the void space between flakes formed by mechanical agglomeration as opposed to actual entrapment of air within the interior of an individual droplet. Thus, even though that patented dispenser uses atomization to make an aerated frozen product, the over-run is determined by the amount of liquid refrigerant adsorbed in the liquid mix before it is fed to the separator. In addition, the extruder is required to further reduce over-run and improve body texture.\nAlso that patented apparatus, like the other dispensers described above, has various internal surfaces and moving parts which must be cleaned in order to keep that dispenser in a sanitary condition. Moreover, there is always a supply of frozen product in the extruder chamber. As described above, this makes it impossible to change quickly from one product flavor to another and, if there is a change without disassembling and cleaning the machine, this results in a mixing of flavors."}
{"text": "Whenever the line or supply voltage is less than the open circuit voltage (OCV) required to operate a gas discharge lamp, the supply voltage magnitude to the lamp must be increased in order to drive the lamp into operation. There must also be some technique to start and restart the lamp, either hot or cold. The required starting voltage is greater than the lamp operating voltage.\nMany different systems have been devised to provide this required operating lamp voltage. The conditions described above, wherein the supply voltage is less than the OCV required for lamp operation, are common because the lowest usable voltage is normally employed for reasons of economy and availability at the application site. One normally uses the highest lumen-per-watt output lamp which is often one of the higher voltage lamps. The lighting system must be consistent with the lighting requirements and must be operable on the available line voltage. If a 120 VAC supply is available, lamps of certain types up to some known wattage level and lumen output can be operated; for the newer, more efficient metal halide lamps and higher wattage lamps, one must arrange for a higher lamp supply voltage such as 240-530 VAC, which may not be available.\nIn these circuits, there are certain basic components, in addition to the lamp itself, which are present, including some form of ballast for voltage transformation and for controlling or limiting the operating current level and lamp power. A semiconductor switching circuit is typically used to step up the source voltage to provide the required lamp ignition and sustaining voltage. A lamp starting circuit is normally present and it is common to switch this starting circuit out of operation, or minimize its influence, after the lamp has entered its normal operation mode.\nStated differently, a lamp operating circuit most often includes a power source, which is normally a low-voltage AC source, some circuit means for controlling the amount of wattage which is delivered to the lamp, and the lamp itself. The circuit usually includes other components for special purposes such as power factor control.\nLamp operating circuits of the prior art have relied upon switching devices such as SCRs, Triacs, transistors or the like to do some of the voltage transformation and control switching, and many of these circuits have included complex and expensive collections of circuits and components. The more components that are used, the more attention that must be paid to the problems associated with heat dissipation and circuit failure rates and life. It is therefore desirable to minimize the number of such components.\nIt is also very desirable, especially in high wattage lamp circuits, to have a high operating power factor for the lamp and the operating circuit. This is sometimes a problem with circuits using large inductive devices, and many circuits of the prior art include capacitive devices to correct the power factor. Switching circuits that are used in lamp operating circuits most often generate a poor power factor and high line harmonics condition."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a contacting and separating mechanism which causes a first member to contact a second member and separates the first member contacting the second member from the second member and an image forming apparatus using the contacting and separating mechanism such as a copying apparatus, a printer, a facsimile machine, and a multifunctional apparatus having the above functions.\n2. Description of the Related Art\nConventionally, in an image forming apparatus such as a copying apparatus and a printer, a contacting and separating mechanism, which causes a member to be contacted such as a secondary transfer roller to contact a contacting member such as an intermediate transfer belt (image carrier) and separates the member contacting the contacting member from the contacting member, has been widely used (see Patent Documents 1 through 3).\nIn an image forming apparatus in Patent Document 1, four photoconductor drums are arrayed to face an intermediate transfer belt. On the four photoconductor drums, a black toner image, a yellow toner image, a magenta toner image, and a cyan toner image are formed, respectively. The four toner images are transferred onto the intermediate transfer belt by being superposed. The superposed toner image on the intermediate transfer belt is transferred onto a recording medium at a position where the intermediate transfer belt contacts a secondary transfer roller.\nIn the image forming apparatus, in order to reduce the wear and deterioration of the intermediate transfer belt and the secondary transfer roller caused by the friction between the intermediate transfer belt and the secondary transfer roller, or in order to remove a recording medium jammed at a position where the intermediate transfer belt contacts the secondary transfer roller, a contacting and separating mechanism is provided which automatically causes the secondary transfer roller to contact the intermediate transfer belt and automatically separates the secondary transfer roller from the intermediate transfer belt.\nSpecifically, in Patent Document 1, by a roller pressure applying lever (pressure applying member) whose rotation supporting point is near the secondary transfer roller, the secondary transfer roller is caused to contact the intermediate transfer belt. When a cam contacting the roller pressure applying lever is rotated to a predetermined angle, the roller pressure applying lever is pushed downward against a spring force which pushes the roller pressure applying lever. With this, the secondary transfer roller is separated from the intermediate transfer belt.\nIn Patent Document 1, in order to replace the secondary transfer roller with a new one, a rotation supporting point for moving the secondary transfer roller by its own weight is provided, in addition to the rotation supporting point of the roller pressure applying lever.\nIn Patent Document 2, a secondary transfer frame pressure applying metal plate and a secondary transfer member (pressure applying member) whose rotation supporting point is near the secondary transfer roller pushes the secondary transfer roller to the intermediate transfer belt. When a cam is rotated to a predetermined angle, which cam is disposed under the secondary transfer roller, and contacts the secondary transfer frame pressure applying metal plate, a spring force pushing the secondary transfer frame pressure applying metal plate and the secondary transfer member is released. With this, the secondary transfer roller is separated from the intermediate transfer belt.\nIn Patent Document 3, a secondary transfer unit (pressure applying member) whose rotation supporting point is separated from the secondary transfer roller causes the secondary transfer roller to contact the intermediate transfer belt. When a cam is rotated to a predetermined angle, which cam is disposed under the secondary transfer roller and contacts the secondary transfer unit, the secondary transfer roller is separated from the intermediate transfer belt.    [Patent Document 1] Japanese Unexamined Patent Publication No. 2004-252258    [Patent Document 2] Japanese Unexamined Patent Publication No. H11-030896    [Patent Document 3] Japanese Unexamined Patent Publication No. 2001-201954\nHowever, in the image forming apparatus, when, for example, a power source of the apparatus is cut off and an automatic contacting and separating mechanism which is moved by a motor is not operated, in order to remove a recording medium jammed at a position where the intermediate transfer belt contacts the secondary transfer roller, or in order to perform maintenance of the intermediate transfer belt and the secondary transfer roller, a manual separating mechanism which manually separates the secondary transfer roller from the intermediate transfer belt must be provided.\nHowever, in a case where the manual separating mechanism is attached to the automatic contacting and separating mechanism, when a manually operable cam (manual cam) contacts the pressure applying member of the automatic contacting and separating mechanism while using the pressure applying member and the rotation supporting point of the automatic contacting and separating mechanism as they are, a great force may be required for the manual operation or pressure deviation may be generated in the width direction of the secondary transfer roller which contacts the intermediate transfer belt.\nThat is, in order to increase the operability of the manual separating mechanism, in a case where the rotation supporting point of the pressure applying member rotating with the secondary transfer roller is disposed near the secondary transfer roller and the manual cam is disposed at a position separated from the secondary transfer roller, when the parallelism between the secondary transfer roller and the rotation supporting point of the pressure applying member is not accurately obtained, the pressure deviation in the width direction of the secondary transfer roller which contacts the intermediate transfer belt becomes great.\nSpecifically, in Patent Documents 1 and 2, since the rotation supporting point of the pressure applying member is near the secondary transfer roller, the pressure deviation in the width direction of the secondary transfer roller which contacts the intermediate transfer belt may be great.\nIn addition, in Patent Document 1, when the secondary transfer roller is replaced with a new one, the rotation supporting point for moving the secondary transfer roller by its own weight is formed, in addition to the rotation supporting point of the pressure applying member. However, the above problem is not solved and the size of the apparatus becomes large.\nIn Patent Document 3, since the rotation supporting point of the pressure applying member is at a position separated from the secondary transfer roller, the pressure deviation of the secondary transfer roller is relatively small at the contacting position with the intermediate transfer belt; however, when the cam is manually moved, a large force is required.\nThe above problems are not limited to the separating mechanism in which the secondary transfer roller is separated from the intermediate transfer belt. That is, the problems are common in all the separating mechanisms in which a member contacting another member is separated from the other member."}
{"text": "Modern inventory systems, such as those in mail order warehouses, supply chain distribution centers, airport luggage systems, and custom-order manufacturing facilities, face significant challenges in responding to requests for inventory items. As inventory systems grow, the challenges of simultaneously completing a large number of packing, storing and other inventory-related tasks become non-trivial. In inventory systems tasked with responding to large numbers of diverse inventory requests, inefficient utilization of system resources, including space, equipment, and manpower, can result in lower throughput, unacceptably long response times, an ever-increasing backlog of unfinished tasks, and, in general, poor system performance. Additionally, expanding or reducing the size or capabilities of many inventory systems requires significant changes to existing infrastructure and equipment. As a result, the cost of incremental changes to capacity or functionality may be prohibitively expensive, limiting the ability of the system to accommodate fluctuations in system throughput."}
{"text": "1. Field of the Invention\nThis invention relates to antennas employing single reflector configurations and, more particularly, to antennas of this character adapted for use simultaneously for a plurality of frequency bands, with the multiple main beams (at different frequencies) coincident, separate or scannable.\n2. Description of the Prior Art\nIt has been proposed in the prior art, for coincident multiple beams, to use a single reflector for a pair of frequency bands from a pair of feeds. This has been achieved by employing dicroic frequency selective surfaces. Prior art systems of this type are disclosed, for example, in T. W. Leonard and L. R. Young, \"Frequency Selective Surfaces,\" AP-S International Symposium, June 1977, pp. 560-563 and A. G. Cha, \"Beam Squint in Large Ground Station Antennas,\" AP-S International Symposium, June 1977, pp. 538-540. Such frequency selective surfaces, however, require high manufacturing tolerances and employ multilayers which increase the losses and cost of the antennas.\nSeparate multiple beams, however, are generated from multiple separate feeds at the focal area of the reflector. These are restricted to minimum beam separation which corresponds to the sizes of the feeding arrangements. The degradation of the beams is a function of the beam separation and the multiplicity of the feeding arrangements. Detailed description of these systems are disclosed, in M. Afifi and P. Foldes, \"Optimum Contiguous Multibeam Antenna Coverage,\" AP-S International Symposium, June 1980, pp. 74-77."}
{"text": "It has been suggested to treat an entire windshield of a vehicle (e.g.—an automobile) with a uniform layer of material that is normally transparent, but fluoresces when it is illuminated with ultraviolet (UV) light, and then project a focused pattern of UV light onto the windshield by using a UV laser to display images on the windshield. The UV laser projects images by altering the direction and intensity of a UV laser beam onto the windshield according to a scanning pattern in order to display both complex indeterminate information, such as a map or other navigation information, and display simple predetermined information such as a symbol indicating that a turn signal is activated. U.S. Pat. No. 6,979,499 issued Dec. 27, 2005 to Walck et al., U.S. Pat. No. 6,986,581 issued Jan. 17, 2006 to Sun et al., and U.S. Pat. No. 7,090,355 issued Aug. 15, 2006 to Liu et al. describe suitable fluorescent materials and application methods, and describe using a UV laser to project images, the entire disclosures of which is hereby incorporated herein by reference. It has been observed that using a UV laser to illuminate such fluorescent material on a windshield may not always provide a sufficiently distinguishable image, such as when bright sunshine is present. Also, when considering providing such a feature on lower priced automobiles, the cost of a UV laser is generally thought to be prohibitive."}
{"text": "A variety of therapeutic articles are used to treat various ailments, injuries and the like. However, a need exists for personal care articles, which can be used in home spa applications. The personal care articles need to be modular and have activatable personal care compositions to provide a variety of therapeutic benefits to a user."}
{"text": "Electronic endoscope devices for a living body to be inserted into a body cavity for internal observation have been developed. In general, such electronic endoscope device is provided with a lens, a solid-state image pickup element (CCD) and the like at a distal end of its insertion portion. An observed image is formed by the solid-state image pickup element through the lens and photoelectrically converted. An electric signal indicating the photoelectrically converted observed image is processed as appropriate by a processor and outputted to a monitor TV, which displays the observed image.\nSize reduction is sought after for from various reasons image pickup devices in which a lens and a solid-state image pickup element are provided at a distal end of an electronic endoscope.\nRegarding a structure to reduce the size of an image pickup device mounted at the distal end of an endoscope, Japanese Patent Application Laid-Open No. 5-115436 discloses size reduction of the image pickup device by mounting parts on a flexible board connected to the solid-state image pickup element and by bending the flexible board in a cylindrical shape so as to enlarge a mounting area."}
{"text": "1. Field of the Invention\nThe present invention relates to a printing apparatus connectable to a plurality of external devices.\n2. Description of the Related Arts\nConventional printing apparatus have been provided for forming copy images based on print data transmitted from a personal computer, word processor or similar data processing device (hereinafter referred to as the \"host\"). Furthermore, printing apparatus connectable to a plurality of hosts have also been provided. The printing apparatus connectable to a plurality of hosts is provided a plurality of connectors for external device input, and print data input from the plurality of hosts to the printing apparatus can be accomplished by connecting signal cables between the aforesaid connectors and the various hosts.\nThe aforesaid type of printing apparatus allows a single printing apparatus to be commonly used as a printing means by a plurality of hosts without providing a selector for selectably switching among the plurality of host devices connected thereto. Such an arrangement is particularly logical from the perspectives of cost and space requirements.\nHeretofore, printing apparatus have been unable to receive print requests from other hosts while printing print data transmitted from a single printing apparatus in accordance with a print request received therefrom. That is, conventional printing apparatus have a disadvantage inasmuch as during the time that the printing apparatus is being used by a single host, that printing apparatus is unavailable for use by other hosts connected thereto."}
{"text": "1. Field of the Invention\nThe subject of the present invention is a method for preparing 19-norsteroid compounds, and the intermediate compounds prepared during the use of this method.\n2. Description of the Art\nOsteoporosis is a bone disease which affects 50 million people worldwide, more particularly women. Its development is linked to age and begins most often after the menopause. This disease is characterized by a reduction in bone density, causes deformations, vertebral compression and in time spontaneous fractures. Osteoporosis therefore represents a serious challenge for public health. The main treatment consists in taking estrogens regularly, which reduces bone loss but which, nevertheless, may be accompanied by certain side effects (bleeding, hot flushes, risk of breast cancer, and the like). A novel series of molecules called SERM (Selective Estrogen Receptor Modulator) allows the treatment of osteoporosis while avoiding some of the side effects (WO98/45316, WO99/67274, WO98/28324, WO99/25725, EP605193, WO02/100880).\nThe subject of the present application is the development of a novel method for preparing a key intermediate (compound of formula I) in the synthesis of certain estrogen derivatives having a dissociated activity.\nAll of the references described herein are incorporated herein by reference in their entirety."}
{"text": "1. Field of the Invention\nThe present invention relates to a vibration motor having brushes with different osculatory tracks, and more particularly to a vibration motor having brushes with different osculatory tracks, which reduces the abrasion of a commutator and has an elongated life span.\n2. Description of the Related Art\nGenerally, communication apparatuses essentially require an incoming-call notification function. The above incoming-call notification function may take the form of sound, such as a melody or bell, or vibration.\nThe vibration is mainly used to prevent the melody or bell from being transmitted to the outside through a speaker, thereby not annoying other persons. In order to generate the vibration, a small-sized vibration motor is driven to transmit its driving force to a case of an apparatus, thereby vibrating the apparatus.\nVibration motors, which are recently applied to portable phones, are divided into coin type vibration motors having a small thickness and bar type vibration motors having a long shape according to shapes of the motors.\nRegardless of the type of vibration motor, the vibration motor uses electromagnetic induction, in which electromagnetic force perpendicular to a magnetic field is generated by flowing current in the magnetic field, to which a conductor is positioned perpendicular.\nThe coin type vibration motor generates vibration in a simple manner such as by rotating a counter weight positioned therein.\nParticularly, the coin type vibration motor has a small thickness and minimizes the number of components of a portable phone, thereby being increasingly used now.\nFIG. 1 is an exploded perspective view of a conventional coin type vibration motor. With reference to FIG. 1, the coin type vibration motor 100 will be described in detail, as follows.\nThe coin type vibration motor 100 comprises a stator assembly 110 serving as a fixed member, and a rotor assembly 120 serving as a rotary member.\nThe stator assembly 110 includes a bracket 111 having a circular flat plate shape, a lower substrate 112 attached to an upper surface of the bracket 111, and a magnet 113 having a disk shape attached to an upper surface of the lower substrate 112.\nThe upper surface of the bracket 111 is sealed by a case 150, and is connected to the case 150 by a central shaft 130.\nThe rotor assembly 120 is rotatably connected to the shaft 130, and includes a bearing 121, coil assemblies 122, a counter weight 123, a commutator 124, an upper substrate 125 and an insulator 126.\nThe upper substrate 125 is a circular printed circuit board, and supplies current having opposite polarities to the coil assemblies 122 through patterns formed on upper and lower surfaces of the upper substrate 125 when power from the commutator 124 attached to the lower surface of the upper substrate 125 is supplied thereto.\nThe commutator 124 is buried under the lower surface of the upper substrate 125 such that a plurality of segments of the commutator 124 form a circle centering on a rotary center and a contact surface of the commutator 124 is exposed to the outside.\nEach of the coil assemblies 122 includes at least one coil so that the coils have the same rotating radius as the magnet 113, and the powers having opposite polarities are supplied to the coil assemblies 122 by the upper substrate 125.\nThe counter weight 123 is made of a material having a high specific gravity, such as tungsten, and is attached on the surface of the upper substrate 125, on which the coil assemblies 122 are positioned, thereby determining the eccentricity of the vibration motor 100.\nThe insulator 126 serves to firmly attach the bearing 121, the coil assemblies 12 and the counter weight 123 to the upper substrate 125, and to insulate the bearing 121, the coil assemblies 12 and the counter weight 123 from the upper substrate 125.\nHere, lower ends of the stator assembly 110 and the rotor assembly 120 are fixed to the lower substrate 112, and upper ends of the stator assembly 110 and the rotor assembly 120 are electrically interconnected by a brush 140 contacting the commutator 124.\nFIG. 2 is a cross-sectional view of a conventional bar type vibration motor. With reference to FIG. 2, the bar type vibration motor 200 will be described in detail, as follows.\nThe bar type vibration motor 200 comprises a stator assembly 210 and a rotor assembly 220.\nThe stator assembly 210 includes a body 211, a cap 212 fixed to one side of the body 211, and a magnet 213. The body 211 has a hollow cylindrical shape, and the magnet 213 is fixed to the inside of the body 211.\nThe rotor assembly 220 includes an eccentric poise 223, a plurality of commutators 224 divisionally fixed to one surface of a fixing body 225, and a plurality of coil assemblies 222 fixed to the fixing body 225.\nThe commutators 224 include a plurality of flat segments on the lower surface of the fixing body 225, and the coil assemblies 222 are electrically connected to the commutators 224 respectively.\nHere, one end of a shaft 230 is connected to the eccentric poise 223, and the other end of the shaft 230 is connected to the fixing body 225. The shaft 230 is rotatably supported by bearings 221 installed in the body 211.\nFurther, a lower substrate (not shown), provided with a pair of brushes 240 fixed thereto, is mounted on the fixing cap 212. The brushes 240 are connected to a lead wire 215 for supplying power therethrough, and contact the commutators 224, thereby applying current.\nThe above-described vibration motors, in which the brushes contact the commutators to supply current, are referred to as “brush type vibration motors”. In order to achieve the commutating action of such brush type vibration motors, it is important to maintain proper stress of the brushes and to make the brushes contact the commutators constantly.\nHereinafter, types of the brushes of the coin-type vibration motor will be described.\nWith reference to FIG. 1, one end of the brush 140 is fixed to the lower substrate 112 by welding, and the uppermost tip of the other end of the brush 140 contacts the commutator 124 so as to apply current. The brush 140 serves to electrically connect the stator assembly 110 and the rotor assembly 120.\nThe brushes 140 are prepared in a pair, which are symmetrically designed for transmitting positive and negative power, and are located on the same osculatory track.\nThe brushes 140 are mainly made of a copper alloy, and the uppermost tips of the brushes 140 contacting the commutators 124 are coated with gold or other precious metal so as to obtain electric conductivity.\nRegardless of which of the above-described types of vibration motors or types of brushes are used, left and right brushes 141 and 142 of the brush unit 140, as shown in FIG. 3, which are symmetrically positioned on the same osculatory track (R), contact the commutators along the same osculatory track (R).\nThat is, since contact portions of the commutators 124 contact the left and right brushes 141 and 142, the contact portions of the commutators are abraded twice per cycle (hereinafter, a track of the brush contacting the rotating surface of the commutator is referred to as an “osculatory track”).\nAccordingly, the contact portions of the commutators are worn out, thus shortening the life span of the vibration motor.\nAccordingly, in order to solve the above problems, a vibration motor having brushes with different osculatory tracks has been required in the art."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to document management. More specifically, embodiments of the present invention relate to a method and system for content dissemination by in-document tagging.\n2. Related Art\nNetwork-based document dissemination is an increasingly critical part of everyday work. However, current dissemination mechanisms, such as email, networked file systems, document management systems, or workflow systems suffer from several limitations.\nIn conventional systems, users can only send information at the level of entire documents. To share selected portions of documents, a user is typically required to manually extract the specific content to be shared, e.g. using cut and paste, and sending it to the recipient as a separate document or in an email message. This manual process not only is cumbersome, but also makes it difficult to reestablish context for the information or to incorporate responses into the original content when the recipient replies with comments or changes to the shared content.\nIn addition, conventional document sharing typically requires special sharing-related interfaces or applications, which are often stand-alone modules and are outside the context in which the information to be shared is actually generated (e.g., a word processor where the document is edited). In addition, most of these interfaces or applications only allow simple distribution of information, whereas specialized functions, such as requests for comments, edits, or approval, are often performed in an ad-hoc fashion using, for example, context provided in an email message."}
{"text": "This invention relates to a composition or mix and method for insulating and inducing slag formation in contact with molten metal, particularly in the manufacture of iron and steel and, more particularly, to an insulating and slag inducing mix which is cast or spread over the surface of the molten metal to form a layer of slag that acts as a refractory thermal insulation and prevents oxidation of the metal by the surrounding atmosphere.\nIn steelmaking and other molten metal processes, a layer of slag on the upper surface of the molten metal prevents both oxidation and excessive heat loss. The slag is formed by the selective oxidation and combination of various impurities (for example, in a charge of blast-furnace iron, scrap iron and other added materials) and causes separation of the material components into desired compositions of steel and slag.\nVarious materials have been used or proposed as additives to form an insulating slag cover for molten metal in steelmaking furnaces. In steelmaking processes, ferrous blast furnace slag has been used as an inexpensive slag inducer to form a layer over the molten metal in basic oxygen or electric arc furnaces. Ferrous slag, however, releases sulfur compounds which cause pollution of the atmosphere. Ferrous slag addition, moreover, creates a thicker than usual layer of slag that increases refractory wear and, in addition, promotes the formation of undesirable sticky or viscous slag. Viscous slag, as has been long known, hinders the rapid transfer of oxygen to the molten metal mass and causes the undesirable build-up of oxygen in the slag in the form of iron oxides.\nMixes of sand and silica have been used as inexpensive fillers but silica can create excessive amounts of slag as well as sticky slag formations and can contaminate the molten metal. Pitch or tar has also been blended with such mixes to provide a less expensive slag inducer. However, the carbon content of pitch and tar gives rise to fire safety and contamination problems. In order to reduce slag formation, aluminum oxide or alimina can be added to the mix. Free or nonfused aluminum oxide, however, has a tendency to plug the furnace orifices and causes undesirable bridging.\nThus, it is apparent that satisfactory slag and slag inducer compositions must have good chemical and thermal stability while in contact with the molten metal and must neither be subject to overly rapid breakdown or excessive slag formation, if contamination of the molten steel as well as air polluting gas formations are to be avoided. Furthermore, as in most industrial processes, it is important that materials selected for inducing slag formation be chosen on the basis of cost effectiveness."}
{"text": "A number of techniques are available for treating cardiovascular disease such as cardiovascular by-pass surgery, coronary angioplasty, laser angioplasty and atherectomy. These techniques are generally applied to by-pass or open lesions in coronary vessels to restore and increase blood flow to the heart muscle. In some patients, the number of lesions are so great, or the location so remote in the patient vasculature that restoring blood flow to the heart muscle is difficult. Percutaneous myocardial revascularization (PMR) has been developed as an alternative to these techniques which are directed at by-passing or removing lesions.\nHeart muscle may be classified as healthy, hibernating and “dead”. Dead tissue is not dead but is scarred, not contracting, and no longer capable of contracting even if it were supplied adequately with blood. Hibernating tissue is not contracting muscle tissue but is capable of contracting, should it be adequately re-supplied with blood. PMR is performed by boring channels directly into the myocardium of the heart.\nPMR was inspired in part by observations that reptilian hearts muscle is supplied primarily by blood perfusing directly from within heart chambers to the heart muscle. This contrasts with the human heart, which is supplied by coronary vessels receiving blood from the aorta. Positive results have been demonstrated in some human patients receiving PMR treatments. These results are believed to be caused in part by blood flowing from within a heart chamber through patent channels formed by PMR to the myocardial tissue. Suitable PMR channels have been burned by laser, cut by mechanical means, and burned by radio frequency current devices. Increased blood flow to the myocardium is also believed to be caused in part by the healing response to wound formation. Specifically, the formation of new blood vessels is believed to occur in response to the newly created wound.\nWhat remains to be provided are improved methods and devices for increasing blood perfusion to the myocardial tissue. What remains to be provided are methods and devices for increasing blood flow to myocardial tissue through controlled formation of channel patterns in the myocardium."}
{"text": "1. Field of Invention.\nThis invention relates generally to a medical ventilator and more particularly to a medical ventilator which is operable in both an invasive and non-invasive ventilatory environment.\n2. Description of the Related Art.\nVentilators are used by patients in various health situations. Typically, these patients have weak physiological attributes that prevent them from breathing unassisted and require invasive ventilation. Invasive ventilatory support generally requires the patient having either a tracheotomy or endotracheal tube disposed into the trachea of the patient. Such treatment generally occurs in hospitals and is administered in acute care situations or post operative recovery situations. Ventilators, such as the Siemens 300, are known to provide invasive ventilatory assistance.\nThe use of invasive ventilatory support has many risks. By positioning an endotracheal tube down the trachea of a patient, the patient is placed at some medical risk. Typically, there are certain physiological attributes in the trachea of a person for preventing bacteria and the like from invading into a person\"\"s lung. However, the endotracheal tube circumvents these natural defense systems. Accordingly, patients may acquire pneumonia by having the bacteria circumvent the natural defense mechanism of the body. Additionally, positioning of an endotracheal tube into a patient subjects the patient to a risk of tracheal abrasion. Overall, the positioning of an endotracheal tube, known as intubation, should be avoided when necessary.\nTypically, when a patient is placed on ventilatory support, most patients are subsequently weaned from the ventilator. Weaning involves manipulating the ventilator from that mode when the ventilator provides the most ventilatory assistance to the mode when the patient is providing most of the breathing. However, during the weaning of the patient from the ventilator, the patient remains intubated continuing further exposing the patient to health risks. Thus, there is a need for a medical ventilator that will enable a patient to be extubated, having the endotracheal tube removed from the trachea, as soon as possible to eliminate health risks.\nAdditionally, there are situations which arise wherein a patient having difficulty breathing is intubated prematurely and connected to a ventilator. This generally occurs since the attending physician lacks a ventilatory device which can provide ventilation in a non-invasive environment for initially determining if the patient is merely having difficulty breathing or truly requires invasive ventilation. This situation arises due to the costs associated with having a separate invasive ventilator and a non-invasive respirator which may be utilized to provide the patient with initial ventilation support. Due to the costs and expenses of having duplicate machines, most hospitals merely have invasive ventilators at their disposal.\nAlso, to provide ventilatory assistance in hospitals, the hospitals generally have specially designed respiratory care facilities having compressed air and oxygen hookups permanently affixed in a specific location. Such fully equipped facilities are expensive and also limit the areas where the ventilator may be accessed. Consequently, some patients who are otherwise healthy but require invasive ventilation are prevented from discharge due to their dependency on the respiratory care facility. With the high cost of hospital stays, some patients who require long term ventilatory utilization may occur exceptional hospitalization charges for the use of such expensive facilities. Thus, it is desirous to enable the ventilatory patient to be discharged to a low-acuity subacute facility or nursing home if the situation is warranted. However, most of these facilities lack the necessary pressurized air and oxygen hookups thus preventing the ventilatory patient from being discharged.\nAlso, since surgical areas and emergency room areas are typically stressful environments, it is desired that ventilators are easy to operate. This is also essential in today\"\"s health care environment since many different type of health care providers are assisting patients. These include respiratory therapists, nurses and physicians. Accordingly, it is desired that an intuitive ventilator exists for both invasive and non-invasive ventilation.\nAccordingly, it is an object of the present invention to provide a ventilator that can be used in both an invasive and non-invasive environment.\nFurthermore, it is an object of the present invention to provide a ventilator which can operate in both an invasive and non-invasive environment and having an operator interface that is simple to use for reducing error in operation;\nAlso, it is an object of the present invention to provide a self contained invasive/non-invasive ventilator having its own source of air and having the ability to mix with an oxygen or other gas source to provide flexibility in providing ventilatory assistance at different physical locations without requiring special respiratory care facilities with preexisting air sources.\nThe above objectives are accomplished according to the present invention by providing a ventilatory system for use in an invasive and non-invasive ventilator environment. The ventilator system includes a gas flow generator for providing a flow of gas to a patient. A conduit delivers the gas flow to the airway of the patient. At least a first valve regulates the delivery of the gas from the gas flow generator to the conduit. A controller controls the delivery of the gas flow to the patient. The system further includes a first set of operational parameters for directing the controller to control the delivery of gas to a patient if the patient is being ventilated in an invasive ventilation mode and a second set of operational parameters for directing the controller to control the delivery of gas to the patient if the patient is being ventilated in a non-invasive ventilation mode. Also a selector is utilized for selecting either the first or second set of parameters to direct the ventilator to provide either invasive or non-invasive ventilatory support to the patient.\nAlso, a unique blower is utilized for providing ventilatory assistance. The blower is a multi-stage centrifugal blower having an air inlet for receiving air from the ambient environment. A first impeller imparts centrifugal force onto the air. A first stator receives the air from the impeller and pressurizes the air. A second impeller subsequently receives the air from the first stator and imparts additional centrifugal force onto the air. A first impeller spacer directs the air from the first stator to the second impeller. A second stator receives the air from the second impeller and further pressurizes the air. A third impeller receives the air from the second stator and further imparts centrifugal force into the air. A second impeller spacer directs the air from the second stator to the third impeller. A blower outlet permits the pressurized air to leave the blower assembly. The overall impeller and stator configuration enables air to be pressurized to at least one hundred and sixty centimeters H2O when exiting the blower outlet.\nAlso, a graphical user interface is utilized in the ventilator for controlling the operation of the ventilator. An activation area is present for displaying a first activation device for a first mode of ventilation and a second activation device for a second mode of ventilation. A selector selects either the first or second activation device. A display area then displays the operational parameters pertaining either to the first or second mode of ventilation depending on the mode selected by the operator. Only those operational parameters relating to the particular ventilation mode selected by the operator are displayed.\nThese and other objects, features, and characteristics of the present invention, and the combination of parts and economies of manufacture, will become more apparent upon consideration of the following detailed description and appended claims with reference to the accompanying drawings, all of which form a part of the specification, wherein like reference numerals designate corresponding parts in the various FIGS."}
{"text": "1. Field of the Invention\nThe present invention relates to a muffler arranged in an exhaust path for reducing exhaust noise.\n2. Description of the Related Art\nFIG. 1 shows a muffler according to a related art. In the muffler 50, exhaust gas passes through an upstream pipe 54 and enters a first chamber 51 through small holes 54a. Also, the exhaust gas exits from an outlet end of the upstream pipe 54 and enters a second chamber 52. The exhaust gas in the first chamber 51 enters a downstream pipe 55 through an inlet end thereof. The exhaust gas in the second chamber 52 enters the downstream pipe 55 through small holes 55a of the downstream pipe 55. The exhaust gas is then discharged into the atmosphere. According to the flow of exhaust gas mentioned above, noise from an engine passes through the upstream pipe 54 and enters the muffler 50. The noise expands when the exhaust gas enters the first chamber 51 through the small holes 54a and contracts when the exhaust gas flows into the downstream pipe 55. Also, the noise expands when the exhaust gas enters the second chamber 52 from the upstream pipe 54 and is muffled when the exhaust gas enters the downstream pipe 55 through the small holes 55a. "}
{"text": "The invention relates to functionalized saccharides and methods of making and using such saccharides.\nThe early diagnosis and treatment of cancer, a family of diseases that affects one in three Americans, continues to present a major challenge (1). While primary tumors can be removed surgically, by the time they are identified, metastatic cells may have spread throughout the body resulting in the lethal growth of tumors in secondary sites. As a result, most treatment protocols involve the additional use of chemotherapeutics, many of which are characterized by severe toxic side effects. These problems have motivated the search for cancer-cell specific agents capable of facilitating the early diagnosis of tumors (i.e., before metastases have developed) and targeting the toxic effects of drugs to cancer cells and away from normal cells. A central feature of these efforts is the identification of cell-surface antigens that are specific to cancer cells. Several such antigens have been identified and, likewise, some progress has been achieved using monoclonal antibodies (mAbs) as delivery agents for drugs (xe2x80x9cimmunotoxinsxe2x80x9d) (2) and diagnostic probes (xe2x80x9cimmunodiagnosticsxe2x80x9d) (3). Unfortunately, the heterogeneity of cancer-associated epitopes has necessitated laborious and cumbersome mAb preparation for each of the myriad of different cancer-associated antigens (4,5). Furthermore, the murine-based mAbs in common usage have proven to be immunogenic in human patients (6) and attempts to xe2x80x9chumanizexe2x80x9d such mAbs have not yet reached maturity (7). A particularly refractory problem is the development of resistant cells that are able to mask or downregulate the targeted antigens (4,8). Thus, alternative approaches are urgently needed for selective cancer cell targeting.\nThe composition of cell surface carbohydrates is dramatically altered in many epithelial and blood-derived cancers. Alterations in cancer-associated cell surface carbohydrate structures can be divided into two broad categories: a) increased levels and inappropriate expression of normal oligosaccharide epitopes and b) expression of novel carbohydrate moieties that do not normally occur on healthy cells. The major antigen families associated with transformed cells are the sialyl Lewis a, sialyl Lewis x, Lewis y, A and Tn-related oligosaccharides (4). A common feature of these antigens is that each contains one or more terminal sialic acid or fucose residues, and the Lewis antigens contain both. The overexpression of sialic acid and fucose residues is highly correlated with increased metastatic potential in several cancers with the greatest impact on human health including gastric, uterine endometrial, colonic, epithelial, pancreatic, bladder, liver, lung, prostate and breast cancers as well as several types of leukemia (9). Consequently, therapeutic and diagnostic strategies that target sialic acid and/or fucose residues may have broad applicability to a variety of cancers.\nThe invention provides methods and compositions for making and using functionalized glycoconjugates. For example, novel and/or unnatural sugars are incorporated into cell associated oligosaccharides using resident pathways of oligosaccharide biosynthesis to expand the informational repertoire of the plasma membrane.\nDisclosed methods for making glycoconjugates include (a) contacting a cell with a first monosaccharide, and (b) incubating the cell under conditions whereby the cell (i) internalizes the first monosaccharide, (ii) biochemically processes the first monosaccharide into a second monosaccharide, (iii) conjugates the second monosaccharide to a carrier to form a glycoconjugate, and (iv) extracellularly expresses the glycoconjugate to form an extracellular glycoconjugate comprising a selectively reactive functional group. In a particular embodiments, the first monosaccharide comprises a chemically reactive functional group such as a ketone, which is incorporated into the second monosaccharide, the glycoconjugate and the extracellular glycoconjugate, the first functional group is N-linked in the first monosaccharide, and the first monosaccharide comprises ManLev.\nThe invention also includes methods for forming a wide variety of products at a cell. The products may provide a label, a binding site, a modulator of cell function such as a drug or toxin, a radiative emission, etc. These methods comprise the steps of making a glycoconjugate according to the invention and then contacting the functional group of the extracellularly expressed glycoconjugate with an agent which selectively reacts with the functional group to form a product. In a particular embodiment, the agent comprises a functional group moiety, such as a hydrazide, which selectively reacts with the functional group of the extracellular glycoconjugate to form a covalent bond, and an effector moiety, such as a drug, which modulates a function of a cell.\nThe subject compositions include cyto-compatible monosaccharides comprising a nitrogen or ether linked functional group, such as a ketone, selectively reactive at a cell surface and compositions and cells comprising such saccharides."}
{"text": "The Health Insurance Portability and Accountability Act of 1996 (HIPAA), the fundamental privacy principles of both the Data Protection Act and the Human Rights Act 1998, along with the American Recovery and Reinvestment Act (ARRA) in February 2009 and Accountable Care Act (ACA) followed by April rulings by the Federal Trade Commission (FTC) included a standard of privacy regarding an individual's right to privacy regarding health care data. In January 2013, a new revision of HIPAA 1996, labeled the HIPAA Omnibus Rule, was issued with increased emphasis on privacy, disclosure of identifiable information and tougher security provisions which comes under the 2009 HITECH Act and the Genetic Information Nondiscrimination act. Under the provisions of HIPAA, ARRA, ACA and the FTC, health information, with few exceptions, can only be shared with the express permission, advance consent, and authorization of the patient (or the patient's legal guardian, as appropriate), and when compromised, electronic notifications (transactions) must be sent, and followed up with electronic audits and risk analysis.\nBy way of example, if a patient is unconscious and has provided advance authorization and consent for a licensed health care provider to securely access and view health-related and protected health information with family, next-of-kin, friends, or others involved, the patient's care and emergency care can be shared when in the best interest of the patient.\nIn Florida, vehicle owners can securely store emergency contact information electronically, including the name and telephone number of at least one person, and link same to their driver's licenses (DL). A law enforcement officer or first responder, if they can locate a driver's license at an accident scene, can contact the Department of Motor Vehicles to obtain emergency contact (ER-Cont) data. If not available and the vehicle occupants are unconscious or otherwise unable to communicate, notifying the family can be a challenge. ER-Cont information is only available to police at a crash scene in the state of Florida.\nNLETS, the National Law Enforcement Telecommunications System, part of Department of Justice, can interface with Department of Motor Vehicle sites across the country and obtain emergency contact information, but only if linked to a vehicle's vehicle identification number (VIN) and with the driver's consent. However, medical data cannot be collected, stored, accessed, or shared via NLETS, which can cause a loss of critical time gaining access to critical healthcare data, such as allergies, blood type, and other medical data. Such data can save lives or improve the quality of life after a life-threatening event. In addition the Government, in 2013, added the Federal ‘First Net’ a telecommunication secure spectrum for Emergency Responders, part of the Department of Commerce.\nAs will be addressed throughout this disclosure, attributes contain information about a subject (known also as an actor). A subject's digital ID has a limited number of identity attributes such as address, age, title or driver's license or trait features that are inherent such as eye color, gender or birth place. A subject can also have acquired associated attributes (lifestyle, purchasing behavior, medical or banking activities) which can change easily whereas trait attributes most likely do not change. Upon being validated and authenticated with a digital ID (public key certificate attributes) and or authenticated attributes in good standing, then a person's (subject) authenticated identity can be enhanced with attributes that originate from an Attribute Certification process where one's Authentication privilege is extended to provide “certified binding attributes' such as access control, secure email, access privileges and associated relationships. As a result of the security and auditing process incorporated into Attribute Certification there is a strong privilege management policy monitoring effort, risk management process and certificate-attribute revocation process. Entities, institutions, exchanges, enterprise servers and the environment (defined as ‘objects’) can also have attributes which are represented by defined characteristics and functions. Attribute certs cannot be used to establish an identity but are used to extend the attributes of one's identity. The forgoing is in concert with NIST guidelines.\nAnonymization and Pseudonymization are specific de-identification processes that file the intent of HIPAA 1996 and the HIPAA omnibus rules of January 2013. Anonymization is the process that removes the identifying characteristics (HIPAA defined) associated with protected health/clinical information and generates a not so unique health data set. The value of such allows a subject/patient to make a part of or subset of their clinical data available for a range of secondary purposes without having to access identifiable clinical information. Such data will be made available on a need to know or on an arranged basis and risk of identity is greatly minimized. The activity is handled through a trusted third party who attests to the validity of the clinical information. Pseudonymization is a specialized class of Anonymization that removes the association and adds an association between a particular set of data characteristics relating to the data subject in addition to adding more pseudonyms. This is a means by which information can be linked together to the same group of persons over time and across multiple data records without revealing the identity of the person and subject data. A trusted third party play's a critical role if there needs to be a re-identification event that is in response to a major public health event. (Activities defined in HIPAA, ACA and HITSP-ONC).\nYet further, there is a need to provide medical help for a patient using a smartphone. By way of example, if the patient is unconscious and has provided advance authorization and consent for a licensed health care provider to securely access and view health-related and protected health information with family, next-of-kin, friends, or others involved, the patient's care and emergency care should be able to be shared when in the best interest of the patient, and in particular during a medical emergency situation where a smartphone provide time access to patient medical information.\nTherefore, it would be beneficial to provide a secure system and method for making both VIN and emergency medical data available on an as-needed basis to licensed emergency medical responders, in order that care be provided in a more efficient, safe, and secure fashion if such data can be voluntarily provided and stored in a secure and separate, non-law-enforcement repository, and linked to the NLETS and the Government's ‘First Net’ telecommunication spectrum for Emergency Responders, both trusted secure infrastructures."}
{"text": "Reduction of the amount of soot particles emitted by diesel engines (diesel particles) during operation can make the use of diesel engines more attractive because of the considerably reduced environmental pollution. The use of soot filters in the exhaust system of diesel engines can lead to the retention of a considerable percentage of the soot particles present in the exhaust gas.\nHowever, the fact that filters that appear to be suitable, especially filters consisting of ceramic threads or ceramic fiber packings, become clogged relatively rapidly, thus losing their filter effect, is problematic.\nTo regenerate such filters, it has been proposed that the soot particles located within the filter material be burned off with a burner. It has also been proposed that an oxidizing agent be sprayed into the filter in order to lower the ignition temperature of the soot.\nSoot filters must be regenerated at certain intervals (determined in terms of hours of operation or the pressure rise of the exhaust gas flowing through). Since the regeneration process takes some time and it cannot be assumed that the engine would be stopped during this period, two diesel particle filters can be arranged in parallel, one of which is always connected in the path of the exhaust gas, while the other is being regenerated and is awaiting use after completion of the regeneration. In an attempt to keep the fuel consumption of an engine as low as possible, the output of the burners during the regeneration of diesel particle filters should also be kept low."}
{"text": "Printing plates are imaged on special opto-mechanical plate setters. The imaged plates undergo a chemical development process which is usually performed by a plate developing device 100 shown in FIG. 1A. The development process uses chemical development material in a liquid form 108. The material contained a developer tank 104 within the plate developer device 100, as is shown in FIG. 1B.\nDuring the filling of the developer tank with developer liquid 108, air 132 may cause carbonization of liquid 108. The carbonization of the developer liquid results in degraded performance of the plate developing process. The intention of this invention is to prevent, or substantially reduce, the carbonization of the developer liquid inside the plate developing machinery."}
{"text": "Fiberoptic illuminators are employed widely in medical, dental and surgical procedures. The illuminator housing encloses a relatively high intensity light source. Metal halide, xenon and halogen light sources are commonly utilized. A fiberoptic cable is connected to a turret or output port mounted on the front panel of the illuminator. Light from the light source is directed through the turret to a focal point located within the turret at or close to the inlet end of the attached fiberoptic cable. The cable conducts the light and delivers it to an attached headlamp or analogous equipment. The light output is then used, as needed, for a particular medical procedure or application.\nThe heat generated by the illuminator at the focal point is typically extremely intense. In conventional illuminators, the turret is composed of a heat conducting metal or metal alloy. As a result, the turret transmits the intense heat at the focal point to the attached fiberoptic cable, as well as to the front panel of the illuminator housing. During a medical or surgical procedure using the illuminator, the doctor, nurse, technician or other medical personnel operating the equipment may have to touch the inlet end of the fiberoptic cable and/or the front panel of the illuminator. Such persons risk being seriously burned by the intense heat that is generated by the illuminator. To date, there are no known illuminator turrets which effectively dissipate this heat and reduce the danger of the user being seriously hurt."}
{"text": "Use of radio frequency bands of the electromagnetic spectrum is regulated by governments in most countries, by allocating specific frequency bands to particular types of uses, such as licensed bands for commercial radio and television broadcasting, cellular telephony, mobile networks such as CDMA2000, WCDMA, HSPA, LTE, and IMT, maritime radio, police, fire, and public safety radio, GPS, radio astronomy, earth stations for satellite communications, and many other uses. Governments also allocate unlicensed bands, for example, for Wireless Regional Area Network (WRAN) broadband access for rural areas and wireless local area networks (WLAN) and wireless personal area networks (WPAN), such as the industrial, scientific, and medical (ISM) band.\nIn the United States, the Federal Communications Commission (FCC) regulates use of the radio spectrum, including radio and television broadcasting. Frequencies are allocated according to a bandplan in which guard bands are assigned between the allocated radio bands to avoid interference between adjacent signals. There are also unassigned frequency bands in the spectrum that either have never been used or have become free as a result of changes in technology. Unassigned or un-used frequencies also appear locally inside the frequency bands, which are otherwise allocated in other locations. The unassigned frequency bands and guard bands are referred to as white spaces.\nTV white space may be broadly defined as broadcast television spectrum that is unused by licensed services. There are at least two categories of TV white space: [1] Dedicated TV white space is a portion of the spectrum that the FCC has reallocated to unlicensed use from previously analog broadcast usage, and [2] Locally unused spectrum by licensed TV broadcasters in a geographic area.\n[1] Dedicated TV white space: In the United States, the FCC has dedicated approximately 400 MHz of white spaces for unlicensed use that became unused after a federally mandated transformation of analog TV broadcasting to digital TV broadcasting. However, the FCC has prohibited unlicensed use of white spaces from interfering with existing licensed uses, including digital TV stations, low power TV stations, cable TV headends, and sites where low power wireless microphones are used. Various proposals have been made for unlicensed use of the white spaces left by the termination of analog TV, for example rural broadband deployment, auxiliary public safety communications, educational and enterprise video conferencing, personal consumer applications, mesh networks, security applications, municipal broadband access, enhanced local coverage and communications, fixed backhaul, and sensor aggregation for smart grid meter reading.\n[2] Locally unused spectrum by licensed TV broadcasters: The FCC has adopted rules to allow unlicensed radio transmitters to operate in the broadcast television spectrum at locations where that spectrum is not being used by licensed broadcasters. The FCC proposes two mechanisms to enable the unlicensed transmitter to discover the available channels: geo-location and database based approach, and spectrum sensing. The use of one of the mechanisms is required for the unlicensed transmitter. The FCC proposed the use of geolocation to establish the location of the unlicensed transmitter and a database of TV bands use by licensed broadcasters organized by their geographic coverage areas, to enable the unlicensed transmitter to know where local TV band white spaces may be available. The FCC proposed the use of spectrum sensors in the unlicensed transmitter to detect the presence of the incumbent, primary TV broadcaster's signal in the local TV band to enable the unlicensed transmitter to immediately relinquish using the band. A primary user in such a local TV band would be an incumbent TV broadcaster licensed to operate in that band, but in those geographic areas where there are no licensed incumbent TV broadcasters in operation, other unlicensed secondary users may make use of that band. There may be also other incumbent users in the TV band, which the secondary users should avoid, such as program making and special events (PMSE) systems.\nIn addition to the United States, other countries are also considering to enable unlicensed, secondary operation in TV band white spaces. The requirements may slightly differ in different countries, e.g. in the United States the maximum transmit power for unlicensed device is defined based on the device type, whereas in Europe location specific maximum transmission power has been considered. In that case the maximum allowed transmission power for an unlicensed device would depend on the device geo-location, i.e. the distance from the primary users. The device characteristics, such as emission mask/ACLR (adjacent channel leakage ratio) may affect the maximum allowed transmission power.\nOther RF spectrum white spaces may be defined as RF spectrum that is locally unused in certain geographic areas, such as for example frequency allocations from maritime radio in landlocked areas remote from the sea. A primary user in such a maritime radio band would be a maritime radio licensed to operate in that band, but in those geographic areas where there are no licensed maritime radios in operation, other unlicensed secondary users may make use of that band. Similarly, locally unused spectrum white spaces may be present in certain geographic locations, such as the frequency allocations from 2.025 GHz to 2.110 GHz for earth stations to transmit to communications satellites, in areas remote from such earth stations. A primary user in such a satellite earth station radio band would be a satellite earth station licensed to operate in that band, but in those geographic areas where there are no satellite earth stations in operation, other unlicensed secondary users may make use of that band. Further, other schemes of secondary use of spectrum, other than unlicensed schemes may exist, such as licensing, regulator defined policies, cognitive principles, or authorized shared access."}
{"text": "This invention relates to optical couplers and more particularly for coupling an arthoscope to a camera such as a television camera.\nAn arthoscope is a device or instrument which allows the practioner to examine the internal structures of joints. This is necessary in performing diagnosis of various joint ailments or diseases.\nIn surgical procedures or for teaching students the practioner desires to couple the arthoscope to a television or other camera while conducting an examination or procedure. In this manner the entire procedure can be viewed on a conventional television receiver.\nThus there are presently available a number of relatively small television cameras which are employed in such procedures. These cameras contain a vidicon tube which converts light images into a signal for application to the television receiver.\nThe television cameras are not constructed uniformly as is known. For example, the vidicon tube or the lens system is never in the same location from camera to camera. Hence a properly focussed camera will operate with an associated arthoscope and another camera will provide an unfocussed picture. The range of focussing of such cameras is limited and this range may not be suitable for operation with various different types of arthoscope. The focal length of the arthoscope from various manufacturers is never the same and the focal length will differ from each manufacturer. In any event, in order to place the image on a television receiver one must have the ability to focus the image prior to application of the same to the television camera.\nIn regard to the arthoscope, these devices are also supplied by various manufacturers. There are of course prior art devices which serve to accomodate or hold various optical instruments to enable one to utilize a camera or other device together with an associated optical instrument. For example, U.S. Pat. No. 4,143,938 entitled MICROSCOPE APPARATUS WITH TELEVISION FILM CAMERA issued on Mar. 13, 1979 to R. E. Feinbloom and assigned to the assignee herein. This Patent shows an optical coupler which has an input port which receives a beam of light from a microscope or beam splitter. One output is coupled to a television camera while another port is coupled to a film camera. The coupler uses a lens asssembly to provide par focal registration at the associated ports and uses a pivotable mirror to select a given output port. In any event, in that Patent, the coupling is afforded by conventional screw threads which coact with the corresponding ports of the coupler.\nIt is an object of the present invention to provide an optical coupler for coupling an arthoscope to a television camera in a simple and reliable manner whereby, the arthoscope can be quickly removed from the coupler. The apparatus also allows one to independently focus the image from the arthoscope prior to directing the image to the television camera."}
{"text": "The present invention is directed to electrical amplifier apparatuses, and especially to electrical amplifier apparatuses adaptively controlling output to accommodate signal clipping to improve recovery from signal clipping.\nThe present invention is generally useful in electrical signal amplifiers. The present invention is particularly advantageously employed in audio signal amplifying apparatuses. Amplification of electrical signals, and especially audio signals, is a challenging area of analog circuit design. Unlike some other linear circuits, audio systems are required to function with high fidelity over a broad dynamic range. While the average power required for most music listening (an example of use of an audio amplifying apparatus) is quite low, there can be peak output levels which are larger than the supply rail—the signal provided for operating the amplifier. When the output signal levels exceed the supply rail, clipping occurs at the outputs of the amplifier. Ideally the amplifier output signal will exactly match the shape of the input signal with a perfect cutoff for levels above the supply voltage. Class-AB power amplifiers typically will have no trouble achieving this kind of performance provided that good design techniques are employed. Class-D amplifiers however, which are rapidly increasing in popularity for audio amplifying applications, present challenging design issues in order to approach such ideal clipping performance.\nA problem that may be experienced during a clipping condition in an amplifier is that, while clipping (i.e., while the output signal levels of the amplifier exceed the supply rail), may be illustrated by way of example and not by way of limitation, in the contest of an exemplary Class D amplifier. The exemplary amplifier receives input signals at an integrator unit. The integrator unit integrates the input signals (using integrating capacitors in a feedback arrangement, which will be described in greater detail hereinafter in connection with FIGS. 1–4. When a clipping condition exists, the amplifier output can no longer maintain proportionality with the input signals received by the amplifier. During this clipping condition, integrating capacitors continue to be charged, so that the output of an integrating unit becomes larger than the reference signal. When the input signal returns to a lower level at which clipping does not occur, the amplifier takes a finite time to reduce the excess charge that has been accumulated in the integrating capacitors during clipping conditions. After the excess charge is reduced, the amplifier can return to normal operations. The delay in reducing excess charge in the integrating capacitors causes distortion during clipping recovery. A user of an audio amplifier system can notice noise occasioned by the distortion that occurs during clipping recovery operation.\nThere is a need for an amplifier apparatus that can reduce time required for clipping recovery operations."}
{"text": "Spread spectrum concepts are gaining increasing importance in the protection of communication links from jamming and detection by unfriendly forces. The most common prior art technique comprises a direct modulation of pseudorandom sequences. A binary code is generated at an information rate which is much greater than that of the data being transmitted. The data is modulo-two added to the binary code, and the resultant sequence is modulated on a suitable carrier frequency, and then it is transmitted. At the receiver, the signal is heterodyned down to baseband to remove the high frequency carrier. A high rate binary sequence which is identical to the transmitter sequence is synchronized to the incoming sequence and modulo-two added to thereby reproduce the lower rate data information back into its original form.\nThe above technique has several disadvantages. First, in the heterodyne procedure, a local oscillator signal of proper frequency is mixed with the incoming signal. To obtain the proper signal at baseband, the phase of the local oscillator must be identical to that of the received signal, thereby necessitating the use of elaborate techniques such as phased lock loops and in-phase and quadrature channels. Second, assuming that the signal has been properly heterodyned, to recover the data, the high-rate digital sequence generated at the receiver must be precisely synchronized with the corresponding sequence in the received waveform. As is known to those skilled in the art, achieving the requisite synchronization is a complex problem."}
{"text": "Third Generation Partnership Project (3GPP) Long Term Evolution (LTE) technology is a mobile broadband wireless communications technology in which transmissions from base stations (referred to as enhanced, or evolved, Node Bs (eNBs)) to mobile stations (referred to as User Equipments (UEs)) are sent using Orthogonal Frequency Division Multiplexing (OFDM). OFDM splits the signal into multiple parallel sub-carriers in frequency. The basic unit of transmission in LTE is a Resource Block (RB), which in its most common configuration consists of 12 subcarriers and seven OFDM symbols (one slot). A unit of one subcarrier and one OFDM symbol is referred to as a Resource Element (RE), as illustrated in FIG. 1. Thus, an RB consists of 84 REs. An LTE radio subframe is composed of two slots in time and multiple RBs in frequency with the number of RBs determining the bandwidth of the system, as illustrated in FIG. 2. Furthermore, the two RBs in a subframe that are adjacent in time are denoted as an RB pair. Currently, LTE supports standard bandwidth sizes of 6, 15, 25, 50, 75, and 100 RB pairs. In the time domain, LTE downlink transmissions are organized into radio frames of 10 milliseconds (ms), each radio frame consisting of ten equally-sized subframes of length Tsubframe=1 ms.\nThe signal transmitted by the eNB in a downlink (the link carrying transmissions from the eNB to the UE) subframe may be transmitted from multiple antennas and the signal may be received at a UE that has multiple antennas. The radio channel distorts the transmitted signals from the multiple antenna ports. In order to demodulate any transmissions on the downlink, a UE relies on Reference Symbols (RSs) that are transmitted on the downlink. These RSs and their position in the time-frequency grid are known to the UE and hence can be used to determine channel estimates by measuring the effect of the radio channel on these symbols. In LTE Release 11 and prior releases of LTE, there are multiple types of RSs. Common Reference Symbols (CRSs) are used for channel estimation during demodulation of control and data messages in addition to synchronization. The CRSs occur once every subframe.\nHeterogeneous, or multi-layer, deployments of LTE networks (referred to as heterogeneous networks) can be used to, e.g., increase capacity and achievable data rates in areas where needed. A heterogeneous network includes a macro cell layer including a number of macro cells served by corresponding eNBs (sometimes referred to as macro eNBs or macro cell eNBs) and a small cell layer including a number of small cells served by corresponding small or small cell eNBs (e.g., Homes eNBs (HeNBs), pico eNBs, micro eNBs, femto eNBs, etc.). The macro eNBs serving the macro cells and the small eNBs serving the small cells have vastly different transmit powers and may be deployed in two main ways. In the first deployment type, the small cell layer and the macro cell layer share the same carrier frequencies, which creates interference between the two layers. In the second deployment type, the small cell layer and the macro cell layer are on separate frequencies.\nMany features of LTE technology benefit from the eNBs in the system being synchronized with each other with respect to transmit timing and frequency. Synchronization of eNBs is typically done using a Global Navigation Satellite System (GNSS), such as Global Positioning System (GPS), or by using network based methods such as IEEE 1588v2. However, such methods may not be available to all eNBs, particularly small cell eNBs in a heterogeneous network. As such, there is a need for systems and methods for transmit timing and frequency synchronization for small cell eNBs in a heterogeneous network."}
{"text": "Many broadcast systems employ one or more headend sites (also known as headends) to receive broadcast signals from multiple sources and/or on multiple channels to facilitate redistribution of the broadcast signals to subscribers and/or other reception sites over a common transmission medium. For example, a typical Internet Protocol Television (IPTV) system employs at least one headend to receive Advanced Television Systems Committee (ATSC) broadcast signals on one or more broadcast channels for redistribution via a digital Internet Protocol (IP) network. The digital format of ATSC broadcasts, as well as the redistribution of their content via a digital IP network, offers the promise of improved signal quality over traditional National Television System Committee (NTSC) analog broadcasts. However, the quality of ATSC broadcast content, like other types of content conveyed by wireless broadcasting systems, can suffer degradation associated with various transmission phenomena or channel conditions, such as co-channel interference, adjacent channel interference and/or weather-induced signal fading."}
{"text": "In diagnosing and treating certain medical conditions, it is often desirable to perform a biopsy, in which a specimen or sample of tissue is removed for pathological examination, tests and analysis. As is known, obtaining a tissue sample by biopsy and the subsequent examination are typically employed in the diagnosis of cancers and other malignant tumors, or to confirm that a suspected lesion or tumor is not malignant. Examination of tissue samples taken by biopsy is of particular significance in the diagnosis and treatment of breast cancer.\nMagnetic resonance imaging (MRI) is often used to locate the site within a breast where a potentially cancerous lesion or tumor is located. Interventional MRI is the magnetic resonance imaging technique (often involving real time imaging) that allows a surgeon to perform MRI-guided tissue biopsy or surgery. One application of interventional MRI is to guide a surgeon during a biopsy or surgical operation on one or both of the breasts of a female patient.\nInterventional MRI procedures typically require a magnetic resonance signal detection coil which has large openings so that the surgeon can have access to the surgical site through the coil with the biopsy device or other surgical devices. The MRI device may also include compression plates to compress the breast. Compression plates contain needle guide holes to allow for proper placement of a needle during a needle biopsy or larger windows configured to hold guide blocks. Guide blocks have block bodies with needle guide holes or with larger openings to allow access for larger biopsy devices. The guide blocks assist the physician with proper placement and stabilization of the biopsy device during a biopsy procedure.\nThe compression plates are secured to the housing of the signal detection coil. Markers, such as Vitamin E Capsules, which are highly visible under MRI device are placed on the compression plates and the position of the markers relative to the suspect tissue is measured using the magnetic resonance images. The proper window of the compression plate is then determined by finding the window in the compression plate which is closest to the desired entry point.\nFor larger biopsy devices the guide blocks are usually cube-shaped and contain about nine circular openings spaced in three rows of three. Due to the configuration of the openings in the guide block and the large size of some biopsy devices it is often difficult to precisely reach the desired tissue location. Present guide blocks for larger biopsy devices do not allow for access to every area of tissue directly behind the block. Thus, there is need in the art for improved guide blocks which allow access to all areas of tissue directly behind the block."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and system for transmitting and receiving an event message. More particularly, the present invention relates to a method and system for transmitting and receiving an event message, which registers phonebook data of users in a server that manages an event message, and which transmits an event message to at least one user having mutually registered phonebook data when an event message transmission is requested.\n2. Description of the Related Art\nRecently, with remarkable advancement of information communication technologies and semiconductor technologies, portable terminals are being widely distributed and used. In particular, recent portable terminals are reaching a mobile convergence phase covering the areas of other terminals beyond traditional unique areas. As a representative example, mobile communication terminals are providing various functions such as a TV viewing function (e.g., a mobile broadcasting such as a Digital Multimedia Broadcasting (DMB) and a Digital Video Broadcasting (DVB)), a music replay function (e.g., an MP3), a photographing function, a data communication function, an Internet connection function, a near field communication function, and the like.\nFurther, as mobile communication services develop and a portable terminal becomes a necessity of citizens, recently, more and more users are informing others of an event (e.g., a first-birthday party, marriage, and advertisement, etc.) using a portable terminal. However, according to the prior art, user had to individually make a phone call or send a text message and an e-mail, etc. using a portable terminal, all of which is relatively inconvenient. Further, in the case in which an advertisement message is transmitted, such an advertisement message is transmitted to unspecified individuals regardless of their thinking. As such, according to the existing message transmission method, it is difficult to expect the effects of advertisement exposure or information provision, and the recipient user may feel unpleasant due to the receipt of spam.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present invention."}
{"text": "The invention relates to covers or closures for containers such as nuclear fuel storage canisters of the type used to store consolidations of spent fuel rods underwater. An understanding of fuel consolidation may be gained from U.S. Pat. No. 4,619,808 issued to Formanek, Oct. 28, 1986 and assigned to Combustion Engineering, Inc., assignee of the instant application.\nIt is important to know that stored spent nuclear fuel rods are in closed canisters which are fully identified as to content and time in residence in a spent fuel pool. To ensure the accuracy of records, indicia are stamped on the canister closure. It is also important to know that the closure has not been tampered with or that access to the contents has not been had to ensure such accuracy. Moreover, it is also important that the closure be inexpensive and capable of providing natural convection flow of the spent fuel pool liquid (water) therethrough. As a practical matter, it is also important that the closure not obstruct canister lift features for remote handling of the canister in the pool and that the tamper-proof closure can be easily positioned in place in the opening of the canister and removed when desired.\nNuclear fuel consolidation canisters are normally provided with a combination lifting and storage cover. These prior art covers are complicated and expensive and one is required for each canister. In using the instant invention, reusable separate lifting covers are employed in handling, but a new and expensive separate lifting cover for each canister is not required, since for storage, only the disposable sheet steel tamper-proof cover of the invention is needed."}
{"text": "A database system may store data across several data volumes of a data area. During operation of the database system, it may occasionally become desirable to delete a data volume from the data area in order to reduce a size of the data area and/or to otherwise reconfigure the data area. Conventional techniques necessitate the database system to enter an offline mode prior to deleting a data volume. Since the database system cannot be read or written in an offline mode, conventional techniques for deleting a data volume are unsuitable for some applications.\nConventional techniques for deleting a data volume pose additional problems if the data volume stores a restart page for the database system. A restart page may include information used to start the database system and reconstruct its underlying structures. Upon start-up, the database system looks for the restart page at a particular location of a particular data volume. The database system will therefore fail to start properly if the data volume including the restart page has been previously deleted.\nSystems are therefore desired for deleting a data volume of a database system while the database system remains in an online operational mode. Systems to perform such deletion and to efficiently move and access a restart page of the deleted data volume are further desired."}
{"text": "1. Field of the Invention\nThe present invention relates to a vibration isolating device having a liquid chamber which is a space defined mainly by a vibration absorbing main body and a flexible membrane and in which a liquid is charged, the liquid chamber being divided into liquid sub-chambers which communicate through a passage, and a method for manufacturing the vibration isolating device.\n2. Description of the Prior Art\nAmong the conventional vibration isolating devices of this kind, a vibration isolating device is well known which comprises an annular vibration absorbing main body made of an elastic material containing a natural rubber having a low loss factor and which is concave at the top thereof and also has attached on the outer circumference thereof a clamping ring to which a top plate is fixed by caulking, a receptacle having an arch-like sectional shape being attached on the inner bottom circumference of the vibration absorbing main body, and a flexible membrane disposed along the bottom surface of the receptacle and having the circumferential edge sandwiched between the receptacle and a bottom plate. There is formed in the top of the receptacle an orifice which provides communication between the liquid chamber formed above the concavity of the vibration absorbing main body and the liquid chamber between the receptacle and flexible membrane (Japanese Utility Model Provisional Publication No. 110344/82 and Japanese Patent Provisional Publication No. 70240/86).\nThe above mentioned conventional vibration isolating device is not of a so-called bush type structure having inner and outer cylinders, and its structure is not suitable for improving the damping effect by lengthening the orifice."}
{"text": "Embodiments of the subject matter disclosed herein generally relate to systems comprising turboexpanders and driven turbomachines and methods for operating the same.\nTurboexpanders are widely used for industrial refrigeration, oil and gas processing and in low temperature processes. In some known applications turboexpanders are used in heat recovery cycles to drive an electric generator. US 2011/0305556 discloses a system and method for power generation including a turboexpander with at least two expansion stages for heat recovery and mechanical power generation to drive an electric generator. In this known application the turboexpander is introduced in a Rankine cycle.\nEP 2400117 discloses the application of a turboexpander-compressor system according to the prior art, wherein the same fluid is processed in the turboexpander and in the compressor. FIG. 1 illustrates the turboexpander-compressor system of the prior art. The system is labeled 200. A turboexpander 210 has a turboexpander impeller 212. The turboexpander 210 receives an inlet gas flow at 214. Inside the turboexpander 210 the gas may expand and thus cause rotation of the turboexpander impeller 212. The expanded gas exits the turboexpander 210 at 216 enabled by a separator (Sep) 220. When the turboexpander-compressor system 200 functions at design conditions, a pressure p1 and a temperature T1 measured respectively by a sensor Sp1 235 and sensor ST1 236 of the inlet gas flow at 214, as well as a pressure p2 and a temperature T2 measured respectively by a sensor Sp2 237 and sensor ST2 238 of the gas flow at the exit side 216 have values close to predetermined values. However, in some situations the turboexpander-compressor system operates in off-design conditions. When off-design conditions occur, the pressure p1 of the incoming gas flow at 214 may be adjusted to become again close to the respective rated value, using, for example, a first set of moveable input guide vanes (IGV1) 218. The first set of moveable input guide vanes 218 are located at an inlet of the turboexpander 210.\nIn the turboexpander-compressor system 200 illustrated in FIG. 1, a compressor 224 has a compressor impeller 226. The compressor 224 receives the gas flow from the turboexpander 210 and delivers a compressed gas flow at the delivery side 228. However, between the turboexpander 210 and the compressor 224, the pressure of the gas flow may be altered due to other process components (e.g., separators, coolers, valves) and pressure losses, so that the gas flow at 216 has pressure p3 when entering the compressor 224 measured by sensor (Sp3) 239.\nThe mechanical work generated by the expansion of the gas in the turboexpander rotates the turboexpander impeller 212. The turboexpander impeller 212 is mounted on the same shaft 230 as the compressor impeller 226. The compressor impeller 226 therefore rotates due to the mechanical work generated during the expansion of the gas in the turboexpander 210. The rotation of the compressor impeller 226 provides energy used to compress the gas in the compressor 224. The mechanical work necessary to rotate the compressor impeller 226 affects the rotating speed u of the shaft 230 measured by speed sensor (Su) 234 and, thereby, indirectly affects the process of expanding the gas inside the turboexpander 210.\nThe turboexpander efficiency is related to a ratio of the rotating speed u of the shaft 230 and the enthalpy drop ΔH across the turboexpander 210. The gas expansion in the turboexpander 210 may be considered approximately an isoentropic process.\nThe characteristic parameters (i.e., p1, T1, p2 and T2) of the gas expansion in the turboexpander 210 and the rotating speed u of the shaft 230 may not vary independently. Therefore, in off-design conditions, in order to maximize the turboexpander efficiency, the pressure p3 of the gas flow at the inlet 216 of the compressor 224 may be controlled, for example, by a second set of moveable inlet guide vanes IGV2 232 provided at the compressor inlet. By modifying the pressure p3 of the gas flow 216 input in the compressor 224, the rotating speed u of the shaft 230 is modified and, therefore, the efficiency of the turboexpander 210 can be maximized.\nA controller 240 receives information regarding the pressure p1 and the temperature T1 of the gas flow at the inlet side 214 of the turboexpander 210, the pressure p3 of the gas flow at the inlet 216 of the compressor 224, and the rotating speed u of the shaft 230, by suitable sensors. The controller 240 may send commands C1 to IGV1 218 in order to adjust the pressure p1 of the gas flow at the turboexpander inlet 214 to be within a predetermined range. Based on monitoring the acquired information, the controller 240 determines when the turboexpander-compressor system 200 functions in off-design conditions. When the controller 240 determines that the turboexpander-compressor system 200 functions in off-design conditions, the controller 240 sends commands C2 to the second set of IGV2 232 to adjust the pressure p3 of the gas input into the compressor in order to maximize a ratio R between the rotating speed u of the shaft 230 and the enthalpy drop ΔH across the turboexpander 210.\nIn this known embodiment, the same controller controls the moveable inlet guide vanes of the turboexpander and the moveable inlet guide vanes of the compressor to optimize the efficiency of the system, based on the assumption that the same fluid is processed in the two turbomachines."}
{"text": "The use of magnetic pulse compression circuits (magnetic modulators) to create high voltage, high current, short duration electrical pulses has been well described (see for example W. S. Melville, \"The use of Saturable Reactors as Discharge Devices for Pulse Generators,\" Radio Section, Paper No. 1034, pp. 185-207, Sep. 15, 1950). Particularly, magnetic modulators have been applied advantageously to drive pulsed gas discharges for excimer lasers and other pulsed gas discharge lasers (see for example Ball et al. U.S. Pat. No. 5,177,754, Jan. 5, 1993.\nPulse propagation delay in a magnetic modulator depends upon the characteristic volt-second product(s) required to saturate the core(s) of the individual magnetic switch element(s). Since the volt-second product is nearly invariant for any individual core, operation at different voltage levels typically results in different pulse propagation or throughput delays relative to a trigger signal from master trigger (a shorter delay at higher voltages and a longer delay at lower voltages). However, in some applications, it is important to maintain a constant propagation delay, regardless of variable operating voltage levels, such that measurement events can be timed relative to a trigger signal from the master trigger. Jitter or variation in timing of the output pulse relative to the trigger signal from the master trigger can also be important in these applications. Additionally, temperature fluctuations can lead to variations in propagation delay and/or jitter in timing of pulses.\nAlthough circuits have been used previously to control pulse propagation delay and jitter in magnetic modulators, they have typically been complex, inflexible, and/or relatively crude and inaccurate in the compensation that they allow (see for example Hill et al. \"Reliable High Repetition Rate Thyratron Grid Driver Used with a Magnetic Modulator,\" 8th IEEE International Pulsed Power Conference, San Diego, Calif., 1991, IEEE Catalog #91CH3052-8).\nA voltage timing compensation circuit in the prior art (see for example Cymer ELS5600 Data DOC. ID: ICLACY00.EPS) compensates for the delay variation at different voltages by adding a low level delay to a trigger signal from a master trigger when the magnetic modulator operates at higher voltages.\nThe prior art voltage timing compensation circuit operates by sampling the voltage of the initial operating stage of a magnetic modulator just prior to master trigger initiation. The sampled voltage is then digitized and used to drive a digital delay generator, which adds to the low level trigger signal a timing delay proportional to the sampled voltage. Thus according to the prior art cited above, delay compensation is linear relative to operating voltage, whereas the actual dependence of pulse propagation delay on voltage is non-linear.\nBecause the characteristic volt-second product of the magnetic switch element(s) is temperature dependent, some prior art implementations comprise a temperature timing compensation circuit to correct for variations in delay caused by fluctuations in operating temperature relative to a nominal design ambient temperature.\nIn accordance with certain prior art embodiments, this temperature timing compensation circuit comprises RC components that synthesize the approximate thermal characteristics of the system.\nWhat is needed in the art is a simple, reliable circuit for timing compensation of a magnetic modulator that more accurately compensates for effects of voltage and/or temperature than does the prior art. Further needed in the art is a circuit for pulse timing control of a magnetic modulator in response to other independent variables, that is simple, accurate, and reliable."}
{"text": "1. Field of the Invention\nThe present invention concerns a device to reduce the electromagnetic field load caused by the presence of a medical intervention apparatus in magnetic resonance examinations.\n2. Description of the Prior Art\nMany medical procedures today are conducted in an interventional manner. Diagnostic or therapy methods that—in contrast to conservative procedures—conduct targeted procedures (interventions) on pathological tissue in order to positively affect the course of the illness are designated as interventional. Invasive or, respectively, minimally invasive procedures are also examples of such methods, wherein medical instruments or implants are brought to the treatment or examination location by various techniques.\nParticularly in the treatment of vascular diseases, guide wires are used that are inserted by a physician into hollow vessels (lumens) of a patient. Later a catheter (for example) is inserted along this wire for diagnosis or therapy. The guide wire itself often has a diameter of less than 0.5 mm. Due to the high mechanical demands, such guide wires have been conventionally produced from metal and metal alloys. Magnetic resonance is a modality particularly suited for imaging of the human body with a high tissue contrast. It would be particularly suitable for observation of intervention apparatuses during their insertion and placement in the human body, but the electromagnetic fields in the two-digit or three-digit megahertz range that are required for imaging are problematical. Therefore, heretofore there have been high risks in the real-time tracking of guide wires and other electrically conductive intervention apparatuses given movement in the human body. This is highly risky because high-frequency currents can form on elongated electrically conductive apparatuses due to injection of the transmission fields. This leads to heating at corresponding powers and, in the disadvantageous case, to tissue burns inside the human body. The currents are even higher at such elongated intervention apparatuses if waveguide resonances occur along the entire length due to the intervention apparatus. If the lengths of the guide wires are also on the order of the wavelength of the magnetic resonance frequencies being used, the risk of endangering a patient is particularly high.\nA prediction about the potential for danger by means of simulation calculations is very difficult since the actual danger is situation-dependent. The length and orientation of the intervention apparatus in the patient, the conductive paths outside the patient, and even properties of the patient, play a role. Finally, the position of the intervention apparatus in relation to the transmission antenna, and the transmission antenna type, are also very influential.\nDue to the high risks from the induced radio-frequency currents, the real-time tracking of intervention apparatuses with the aid of magnetic resonance has previously been realized with the transmission power being drastically limited. Another solution is the construction of guide wires or catheters from non-conductive materials. However, all desired material properties cannot be ensured with these materials; such intervention apparatuses are additionally very expensive and complicated."}
{"text": "1. Field of the Invention\nThe present invention relates to an actuator and a blade drive device for a camera.\n2. Description of the Related Art\nConventionally, there has been known an actuator used for driving a shutter blade employed in a camera, and composed of a rotor, a stator, a coil for exciting the stator, and an outputting member for transmitting the rotational movement of the rotor to the shutter blade. Japanese Unexamined Patent Application Publication No. 2004-309531 discloses a device in which a rotor and a drive pin, serving as an outputting member, are integrally formed.\nHowever, in the case where the rotor and the drive pin are integrally formed, it is difficult to change the positional relationships between the rotor and the drive pin whereas it is effective for downsizing. Therefore, it is difficult to change an angular position of the drive pin relative to the rotor according to a design variation."}
{"text": "The present invention relates to amplifiers and more particular to amplifiers employing shunt feedback.\nThe traditional series shunt feedback architecture illustrated in FIG. 1 has been widely hailed as the preferred architecture for Low Noise Amplifiers (LNAs) requiring a well defined input impedance. However combining series shunt feedback with Gilbert\"\"s current steering techniques for providing an Automatic Gain Control (AGC) function presents certain difficulties.\nThe following equations describe how the input impedance Rin of the traditional series shunt feedback amplifier is set by a negative feedback amplifier:\nRf=Rs*(1+A)\nRin=Rf/(1+A)\nwhere Rf is the effective feedback impedance, Rs is the source impedance, and |A| is the magnitude of the gain of the amplifier. It will be appreciated that as long as Rf remains constant, then so does the input impedance Rin. However, if the amplifier has a varying gain, then the input impedance Rin will also change.\nA potential solution to this problem is offered by the architecture of FIG. 2, where the input impedance of an amplifier arrangement is set by an input stage comprising an amplifier which has a fixed gain. The output of the input stage feeds a variable gain amplifier stage. This solution is not optimal however, since the input stage has a fixed gain and the input signal to the variable gain amplifier stage is therefore higher than the input signal to the amplifier arrangement as a whole. This creates a problem in signal handling requirements and may result in a power in-efficient solution.\nIn view of the disadvantages noted above, existing amplifier arrangements may not meet the strict requirements of modem electronic systems, such as interactive cable television equipment, in which a range of high performance RF products are required ranging from mixer oscillator circuits to LNAs with AGC and loop through circuits.\nIt is an object of the present invention to provide an amplifier arrangement having a stable input impedance. It is a further object of the present invention to provide a low noise amplifier which is power efficient. These and other objects are achieved at least in part by an amplifier arrangement in which the amplifier stages thereof are operated in current mode.\nAccording to a first aspect of the present invention there is provided an AC voltage amplifier arrangement comprising:\nan input stage having an input, an amplifier connected in series shunt feedback mode, and an output; and\na second stage having an input coupled to the output of said input stage, and an output,\nwherein said input and second stages are arranged to operate in current mode.\nPreferably, said second stage is a variable gain amplification stage.\nPreferably, the voltage amplifier arrangement is arranged to amplify radio frequency voltage signals.\nPreferably, said amplifier of the input stage is a transconductance amplifier. More preferably, the transconductance amplifier comprises a bipolar junction transistor configured as a common emitter stage, or a differential pair of bipolar junction transistors configured as common emitter stages.\nPreferably, the arrangement comprises a unity gain buffer amplifier coupled in the series shunt feedback loop of the input stage.\nPreferably, said second stage comprises a current steering arrangement arranged in use to steer a proportion of the current, provided to the input of the second stage, to an output load, whereby an output voltage is developed across said load and is available at said output of the second stage. The current steering arrangement may be a Gilbert steering stage. More preferably, said current steering arrangement comprises a pair of bipolar junction transistors having their emitters coupled together and to said output of the input stage. The current steering arrangement may comprise two pairs of bipolar junction transistors in a differential arrangement, coupled respectively to differential outputs of the input stage.\nThe second stage may be arranged to be controlled by an Automatic Gain Control (AGC) signal in order to maintain the output of the stage at a desired level. More preferably, where the stage comprises a pair (or pairs) of bipolar junction transistors, said AGC signal is applied to the bases of the transistors.\nThe second stage may comprise a mixer or a buffer.\nAccording to a second aspect of the present invention there is provided a method of amplifying an AC voltage signal, the method comprising:\nconverting the voltage signal into a current signal;\namplifying said current signal to a desired level; and\nconverting the amplified current signal into an amplified voltage signal.\nPreferably, the amplification factor by which the current signal is amplified is less than one.\nPreferably, amplification of the current signal comprises splitting the current signal into at least two parts, with one of said parts being converted into the amplified voltage signal."}
{"text": "In automatic transmissions of modern motor vehicles, the transmission control modules are in many cases arranged within the automatic transmission. In order to electrically conductively connect these transmission control modules to electrical connection lines for voltage supply, to signal lines and to bus lines, electrical lines are guided from the transmission control module to an electrical plug-type connector, which is arranged in a through-hole in a transmission housing surrounding the automatic transmission. The transmission housing delimits an interior space filled with automatic transmission oil from an exterior space surrounding the automatic transmission. The electrical plug-type connector can be electrically contacted from the exterior space by means of an electrical mating plug-type connector. The electrical plug-type connector arranged in the through-hole is formed in such a way that no fluid can pass from the interior space to the exterior space. In this case, the plug housing is sealed on its outer wall with respect to the transmission housing, generally by means of an elastomer ring seal. Electrical plug-type contacts are generally arranged in the plug housing of the electrical plug-type connector and protrude from the exterior space into the interior space. These plug-type contacts are sealed with respect to the plug housing fabricated from plastic. This can be achieved for example in that the plug-type contacts have specifically shaped surface profiles, such as a fir-tree profile, in the region with which the plug-type contacts are connected to the plug housing. It is also known to form the seal using O-rings, which are fitted to each plug-type contact. Furthermore, seals by casting, for example with polyurethane, epoxy resin, silicones or similar materials, are known.\nIt has been found however that, after a few temperature changes within the automatic transmission, wherein the temperature fluctuations may be between −40° C. and +150° C., the adhesion or press-fit between the plug-type contact and the plug housing can loosen. Casting compounds are also susceptible to crack formations or are not sufficiently resistant to automatic transmission oils. This leads to leaks in so far as fluid can now pass from the interior space along the plug-type contacts through the plug housing into the exterior space."}
{"text": "The present invention relates to metal coordination complexes having constrained geometry. The invention also relates to certain novel addition polymerization catalysts comprising such metal complexes having constrained geometry. Furthermore, the invention relates to methods for the polymerization of addition polymerizable monomers and to the resulting polymers.\nBecause of the unique exposure of the active metal site of the metal coordination complexes having constrained geometry, catalysts resulting therefrom have unique properties. Under certain conditions, the catalysts of the invention are capable of preparing novel olefin polymers having previously unknown properties due to their unique facile abilities to polymerize xcex1-olefins, diolefins, hindered vinylidene aliphatic monomers, vinylidene aromatic monomers and mixtures thereof.\nNumerous metal coordination complexes are known in the art including such complexes involving monocyclopentadienyl groups and substituted monocyclopentadienyl groups. The present metal coordination complexes differ from those previously known in the art due to the fact that the metal is bound to a delocalized substituted xcfx80-bonded moiety in a manner so as to induce a constrained geometry about the metal. Preferably the metal is bound to a cyclopentadienyl, substituted cyclopentadienyl or similar group by both a xcex75-bond and a bridging linkage including other ligands of the metal. The complexes also preferably include metals having useful catalytic properties.\nAlso previously known in the art are transition metal coordination complexes known as tucked complexes. Such complexes are described in Organometallics 6, 232-241 (1987).\nIn U.S. Ser. No. 8,800, filed Jan. 30, 1987 (published in equivalent form as EP 277,004) there are disclosed certain bis(cyclopentadienyl) metal compounds formed by reacting a bis(cyclopentadienyl) metal complex with salts of Bronsted acids containing a noncoordinating compatible anion. The reference discloses the fact that such complexes are usefully employed as catalysts in the polymerization of olefins. For the teachings contained therein U.S. Ser. No. 8,800 and EP 277,004 are herein incorporated in their entirety by reference thereto.\nPrevious attempts to prepare copolymers of vinylidene aromatic monomers and xcex1-olefins, in particular copolymers of styrene and ethylene, have either failed to obtain substantial incorporation of the vinylidene aromatic monomer or else have achieved polymers of low molecular weight. In Polymer Bulletin, 20, 237-241 (1988) there is disclosed a random copolymer of styrene and ethylene containing 1 mole percent styrene incorporated therein. The reported polymer yield was 8.3xc3x9710xe2x88x924 grams of polymer per micromole titanium employed.\nIt has now been discovered that previously known addition polymerization catalysts are incapable of high activity and polymerization of numerous monomers because they lack constrained geometry.\nIt would be desirable if there were provided novel complexes of groups 3 (other than scandium), 4-10 and the lanthamides having constrained geometry.\nAdditionally it would be desirable if there were provided novel catalysts for addition polymerizations comprising novel complexes of groups 3 (other than scandium), 4-10 and the lanthamides having constrained geometry.\nFurthermore, it would desirable if there were provided a process for the preparation of polymers of addition polymerizable monomers using novel catalysts comprising complexes of groups 3 (other than scandium), 4-10 and the lanthamides having constrained geometry.\nFinally, it would be desirable if there were provided novel polymers of addition polymerizable monomers that may be prepared by an addition polymerization process using catalysts comprising novel complexes of groups 3 (other than scandium), 4-10 and the lanthamides having constrained geometry.\nIn one aspect the present invention relates to a metal coordination complex having constrained geometry. More particularly it relates to such coordination complexes that are usefully employed in combination with activating cocatalyst compounds or mixtures of compounds to form a catalytic system usefully employed in the polymerization of addition polymerizable monomers, especially ethylenically unsaturated monomers.\nIn another aspect the present invention relates to a process for preparing certain components of the above metal coordination complexes having constrained geometry and to the precursor compounds necessary therefor.\nIn yet another aspect the present invention relates to a process for preparing addition polymers, especially homopolymers and copolymers of olefins, diolefins, hindered aliphatic vinyl monomers, vinylidene aromatic monomers and mixtures of the foregoing and to the resulting polymer products.\nAccording to the present invention there is provided a metal coordination complex comprising a metal of group 3 (other than scandium), 4-10 or the lanthamide series of the periodic table of the elements and a delocalized _-bonded moiety substituted with a constrain-inducing moiety, said complex having a constrained geometry about the metal atom such that the angle at the metal between the centroid of the delocalized, substituted xcfx80-bonded moiety and the center of at least one remaining substituent is less than such angle in a similar complex containing a similar xcfx80-bonded moiety lacking in such constrain-inducing substituent, and provided further that for such complexes comprising more than one delocalized, substituted xcfx80-bonded moiety, only one thereof for each metal atom of the complex is a cyclic, delocalized, substituted xcfx80-bonded moiety.\nIn addition there is provided a metal coordination complex corresponding to the formula: \nwherein:\nM is a metal of group 3 (other than scandium), 4-10, or the lanthamide series of the periodic table of the elements;\nCp* is a cyclopentadienyl or substituted cyclopentadienyl group bound in an xcex75 bonding mode to M;\nZ is a moiety comprising boron, or a member of group 14 of the periodic table of the elements, and optionally sulfur or oxygen, said moiety having up to 20 non-hydrogen atoms, and optionally Cp* and Z together form a fused ring system;\nX independently each occurrence is an anionic ligand group or neutral Lewis base ligand group having up to 30 non-hydrogen atoms;\nn is 0, 1, 2, 3, or 4 depending on the valence of M; and\nY is an anionic or nonanionic ligand group bonded to Z and M comprising nitrogen, phosphorus, oxygen or sulfur and having up to 20 non-hydrogen atoms, optionally Y and Z together form a fused ring system.\nThere is also provided according to the present invention a process for preparing a metal coordination complex corresponding to the foregoing formula I comprising the steps of:\nA) contacting a metal compound of the formula MXn+2 or a coordinated adduct thereof with a dianionic salt compound corresponding to the formula:\n(L+x)y(Cp*-Z-Y)xe2x88x922xe2x80x83xe2x80x83(II) \nor \n((LXxe2x80x3)+x)y(Cp*-Z-Y)xe2x88x922xe2x80x83xe2x80x83(III) \nwherein:\nL is a metal of group 1 or 2 of the periodic table of the elements,\nXxe2x80x3 is fluoro, chloro, bromo, or iodo,\nx and y are either 1 or 2 and the product of x and y equals 2, and\nM, X, Cp, R, and Y are as previously defined; and\nB) recovering the resulting product.\nFurther there is provided a process for preparing a metal coordination complex corresponding to the foregoing formula I comprising the steps of:\nA) contacting a metal compound of the formula MXn+1 or a coordinated adduct thereof with a dianionic salt compound corresponding to the formulas II or III;\nB) oxidizing the metal to a higher oxidation state by contacting the reaction product of step A) with a noninterfering oxidizing agent; and\nC) recovering the resulting product.\nThere is also provided a catalyst useful in addition polymerizations comprising the following components:\na) a metal coordination complex comprising a metal of group 3 (other than scandium), 4-10 or the lanthamide series of the periodic table of the elements and a delocalized xcfx80-bonded moiety substituted with a constrain-inducing moiety, said complex having a constrained geometry about the metal atom such that the angle at the metal between the centroid of the delocalized, substituted xcfx80-bonded moiety and the center of at least one remaining substituent is less than such angle in a similar complex containing a similar xcfx80-bonded moiety lacking in such constrain-inducing substituent, and provided further that for such complexes comprising more than one delocalized, substituted xcfx80-bonded moiety, only one thereof for each metal atom of the complex is a cyclic, delocalized, substituted xcfx80-bonded moiety; and\nb) an activating cocatalyst.\nIn a further embodiment of the present invention there is provided a catalyst useful in addition polymerizations comprising the following components:\na) a metal coordination complex corresponding to the formula I, and\nb) an activating cocatalyst.\nEven further according to the present invention there is provided a polymerization process comprising contacting one or more addition polymerizable monomers under addition polymerization conditions with a catalyst comprising:\na) a metal coordination complex comprising a metal of group 3 (other than scandium), 4-10 or the lanthamide series of the periodic table of the elements and a delocalized xcfx80-bonded moiety substituted with a constrain-inducing moiety, said complex having a constrained geometry about the metal atom such that the angle at the metal between the centroid of the delocalized, substituted xcfx80-bonded moiety and the center of at least one remaining substituent is less than such angle in a similar complex containing a similar xcfx80-bonded moiety lacking in such constrain-inducing substituent, and provided further that for such complexes comprising more than one delocalized, substituted xcfx80-bonded moiety, only one thereof for each metal atom of the complex is a cyclic, delocalized, substituted xcfx80-bonded moiety; and\nb) an activating cocatalyst.\nFurther according to the present invention there is provided an addition polymerization process comprising the steps of;\nA) contacting a mixture comprising one or more addition polymerizable monomers under polymerization conditions in the presence of a catalyst comprising a metal coordination complex corresponding to the formula I and an activating cocatalyst; and\nB) recovering the resulting polymer.\nFurther still according to the present invention there is provided a polymer comprising in interpolymerized form one or more addition polymerizable monomers prepared by contacting an addition polymerizable monomer or mixture thereof under addition polymerization conditions with a catalyst comprising:\na) a metal coordination complex comprising a metal of group 3 (other than scandium), 4-10 or the lanthamide series of the periodic table of the elements and a delocalized xcfx80-bonded moiety substituted with a constrain-inducing moiety, said complex having a constrained geometry about the metal atom such that the angle at, the metal between the centroid of the delocalized, substituted xcfx80-bonded moiety and the center of at least one remaining substituent is less than such angle in a similar complex containing a similar xcfx80-bonded moiety lacking in such constrain-inducing substituent, and provided further that for such complexes comprising more than one delocalized, substituted xcfx80-bonded moiety, only one thereof for each metal atom of the complex is a cyclic, delocalized, substituted xcfx80-bonded moiety; and\nb) an activating cocatalyst.\nIn a still further embodiment of the present invention there is provided a polymer comprising in polymerized form one or more addition polymerizable monomers prepared by contacting an addition polymerizable monomer or mixture thereof under addition polymerization conditions with a catalyst comprising the following components:\na) a metal coordination complex corresponding to the formula I and an activating cocatalyst.\nIn still further embodiments there are provided ElPE polymers which are highly elastic, interpolymers of ethylene and one or more olefins other than ethylene.\nIn addition there are provided pseudo-random interpolymers of an xcex1-olefin, particularly ethylene and a vinylidene aromatic monomer, a hindered aliphatic vinylidene monomer, or a mixture thereof.\nThe complexes of the invention are usefully employed as catalysts for addition polymerization processes to prepare polymers that are useful as molded articles, films for packaging applications, and foams for cushioning applications; and in the modification of synthetic and naturally occuring resins. The complexes may also be used as catalysts for hydrogenations, catalytic cracking processes, and in other industrial applications."}
{"text": "Optical gratings are well known in the art and are used to disperse optical spectra spatially. Such gratings are commonly used in optical spectrometers to analyze the spectral composition of an optical beam. There is always a tradeoff between the length of an optical spectrometer and its resolution. Thus, if a higher wavelength resolution is required, the length required is also longer. Consider an example of a typical 1-meter long grating spectrometer in the market, which has a typical wavelength resolution of about ΔλSM=0.1 nm at a center operating wavelength λSM=1000 nm or ΔλSM/λSM=10−4. The dimensionless quantity for the length of the spectrometer LSM is LSM/λSM and LSM/λSM=106 in this example. The dimensionless product of the relative resolution ΔλSM/λSM and the relative physical size LSM/λSM of the spectrometer is dependent on the design of the spectrometer.\nResolution Vs Size/Area Factor\nMore precisely, consider a curved grating spectrometer device 1000 shown in FIG. 1 with a curved grating CG 1010 and a wave or beam propagating region between the input slit SLI1 1201 and the grating or the grating and the output slit SLO1 1401 called grating-propagating region GPR 1020. If the spectrometer has a width WSM 1002 and a length LSM 1003, then the dimensionless product of interest is:RSSM=(ΔλSM/λSM)*[(WSM/λSM)*(LSM/λSM)]0.5,  (1)where “*” denotes numerical multiplication. This factor (RSSM) is referred to here as the “resolution vs size” factor of the spectrometer. The square of RSSM factor is called “resolution vs area” factor:RASM=(RSSM)2.  (2)The factors RASM and RSSM basically measure the compactness of a spectrometer for a given spectral resolution power. In this example, if WSM=LSM=1 m, then RSSM=100 and RASM=104. The smaller the RSSM and RASM value, the more compact is the spectrometer.\nOnly a few conventional spectrometers have RASM factor less than about 100 (or RSSM<10). This is primarily because of the various limitations in the prior arts (as will be described below). In many photonic integrated circuit, electronic-photonic integrated circuit, and fiber-optic applications, RASM factors of smaller than about RASM=1 down to RASM=0.01 are desirable. These cannot be achieved via the methods in the prior arts.\nFurthermore, for many applications in photonic integrated circuits, electronic-photonic integrated circuits, and fiber-optics, the background extinction of the spectral power outside the wavelength range that is supposed to be detected is important, which measures the relative optical power from outside the detection wavelength range that is scattered into the output slit or photodetector of the spectrometer.\nSpectral Output Power Efficiency\nA parameter that measures such unwanted background scattering is the “adjacent-channel power extinction ratio”. As shown by FIGS. 1 and 1B, assuming the input light at an input slit labeled as slit “SLI1” 1201 for which “I1” labels “input” “1” has an optical spectrum with optical power spectral density PSI1(λ) 1131 at wavelength λ. Let λI1-O1 1321 be the center spectral wavelength that goes from the spectrometer's input slit SLI1 to the spectrometer's output slit labeled as output slit “SLO1” 1401. Let the optical power spectrum entering or being detected at output slit SLO1 1401 be PSI1-O1(λ) 1331. PSI1-O1(λ) 1331 is then given in terms of the spectral density of the input beam PSI1(λ) by:PSI1-O1(λ)=ηeffI1-O1(λ)*PSI1(λ).  (3)The factor ηeffI1-O1(λM) is the efficiency of passing or detecting the power of the input beam from input slit SLI1 at the output slit SLO1 at wavelength λM, assuming that the input beam is basically a monochromatic light source at λM (e.g. a narrow-bandwidth laser light). The factor ηeffI1-O1(λM) is called the “spectral power output efficiency”.\nAs shown in FIG. 1B, the power at input slit SL 1201 over a small spectral bandwidth Δλ centered at wavelength λA (small comparing to the spectral bandwidth of PSI1-O1(λ) at λA or more precisely, small enough so that PSI1-O1(λ) at λA does not change much over the wavelength bandwidth Δλ), denoted as PI1(λA; Δλ) 1131DL, is given by:PI1(λA;Δλ)=PSI1(λA)*Δλ.  (4A)Let the total optical power entering or being detected at output slit SLO1 1401 over a small spectral bandwidth Δλ centered at wavelength λA for the beam from input slit SLI1 be PI1-O1(λA; Δλ) 1331DL. PI1-O1(λA; Δλ) is then given by PSI1-O1(λA)*Δλ=ηeffI1-O1(λ)*PSI1(λ)Δλ according to Eq. (3). It is also given by ηeffI1-O1(λA)*(λA; Δλ). PI1-O1(λA; Δλ) is thus related to the spectral density of the input beam PSI1(λA) by:PI1-O1(λA;Δλ)=PSI1-O1(λA)*Δλ=ηeffI1-O1(λA)*PSI1(λA)*Δλ.  (4B)In the situation that Δλ is large, Eqs. (4A) and (4B) should be more precisely converted to an integration of PSI1-O1(λ) with respect to wavelength λ over wavelength bandwidth Δλ centered at wavelength λ=λA given by:\n                                              ⁢                                            P                              I                ⁢                                                                  ⁢                1                                      ⁡                          (                                                λ                  A                                ;                Δλ                            )                                =                                    ∫                                                λ                  A                                -                                  Δλ                  2                                                                              λ                  A                                +                                  Δλ                  2                                                      ⁢                                                            PS                                      I                    ⁢                    1                                                  ⁡                                  (                  λ                  )                                            ⁢                                                          ⁢              d              ⁢                                                          ⁢              λ                                                          (                  5          ⁢          A                )                                                      P                                          I                ⁢                1                            -                              O                ⁢                1                                              ⁡                      (                                          λ                A                            ;              Δλ                        )                          =                                            ∫                                                λ                  A                                -                                  Δλ                  2                                                                              λ                  A                                +                                  Δλ                  2                                                      ⁢                                                            PS                                                            I                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      ⁡                                  (                  λ                  )                                            ⁢                                                          ⁢                              ⅆ                λ                                              =                                    ∫                                                λ                  A                                -                                  Δλ                  2                                                                              λ                  A                                +                                  Δλ                  2                                                      ⁢                                                            η                                                            effI                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      ⁡                                  (                  λ                  )                                            *                                                PS                                      I                    ⁢                    1                                                  ⁡                                  (                  λ                  )                                            ⁢                                                          ⁢              d              ⁢                                                          ⁢              λ                                                          (                  5          ⁢          B                )            Total Power and Power Spectrum\nTo define the “adjacent-channel power extinction ratio”, let the spectral resolution bandwidth for a beam at wavelength λI1-O1 1321 from input slit SLI1 1201 to output slit SLO1 1401 be ΔλRes-I1-O1 1381Res, as shown in FIG. 1B. Assuming the input power spectrum at an input slit of the spectrometer PSI1(λ) 1131 is approximately a constant for the wavelength λ around λ=λI1-O1 1321. Let the optical power from the input slit SLI1 1201 entering or being detected at an output slit of the spectrometer SLO1 1401 over the spectral bandwidth given by the spectral resolution bandwidth ΔλRes-I1-O1 1381Res be PI1-O1(λI1-O1; ΔλRes-I1-O1), where:\n                                          P                                          I                ⁢                1                            -                              O                ⁢                1                                              =                                                    P                                                      I                    ⁢                    1                                    -                                      O                    ⁢                    1                                                              ⁡                              (                                                      λ                                                                  I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              ;                                      Δλ                                          Res                      -                                              I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                                            )                                      =                                          ∫                                                      λ                                                                  I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              -                                                            Δλ                                              Res                        -                                                  I                          ⁢                          1                                                -                                                  O                          ⁢                          1                                                                                      2                                                                                        λ                                                                  I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              +                                                            Δλ                                              Res                        -                                                  I                          ⁢                          1                                                -                                                  O                          ⁢                          1                                                                                      2                                                              ⁢                                                                    η                                                                  effI                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              ⁡                                      (                    λ                    )                                                  *                                                      PS                                          I                      ⁢                      1                                                        ⁡                                      (                    λ                    )                                                  ⁢                                                                  ⁢                d                ⁢                                                                  ⁢                λ                                                    ,                            (                  6          ⁢          A                )            Then at an adjacent wavelength that is separated from λI1-O1 1321 by “one spectral resolution bandwidth” ΔλRes-I1-O1 1381Res away, the power detected over the same spectral resolution bandwidth ΔλRes-I1-O1 1381Res given by:\n                                                        P                              I1                -                                  O                  ⁢                  1                                                      ⁡                          (                                                                    λ                                                                  I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              +                                      Δλ                                          Res                      -                                              I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                                            ;                                  Δλ                                      Res                    -                                          I                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                                  )                                =                                    ∫                                                λ                                                            I                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      +                                  Δλ                                      Res                    -                                          I                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      -                                                      Δλ                                          Res                      -                                              I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              2                                                                              λ                                                            I                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      +                                  Δλ                                      Res                    -                                          I                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      +                                                      ΔλΔλ                                          Res                      -                                              I                        ⁢                        1                                            -                                              O                        ⁢                        1                                                                              2                                                      ⁢                                                            η                                                            effI                      ⁢                      1                                        -                                          O                      ⁢                      1                                                                      ⁡                                  (                  λ                  )                                            *                                                PS                                      I                    ⁢                    1                                                  ⁡                                  (                  λ                  )                                            ⁢                                                          ⁢              d              ⁢                                                          ⁢              λ                                      ,                            (                  6          ⁢          B                )            should ideally be zero at the output slit (as the spectral power output efficiency ηeffI1-O1(λ) shall be approximately zero when λ=λI1-O1+ΔλRes-I1-O1 is not at the supposed output detection/passing wavelength λI1-O1 1321 of output slot SLO1 and is at one resolution bandwidth ΔλRes-I1-O1 1381Res away from λI1-O1 1321).Adjacent-Channel Extinction Ratio\nThe ratio of the power at λI1-O1 1321 and the power at λI1-O1+ΔλRes-I1-O1 given by:ηace(λI1-O1)=PI1-O1(λI1-O1;ΔλRes-I1-O1)/PI1-O1(λI1-O1+ΔλRes-I1-O1;ΔλRes-I1-O1),  (7)is called the “adjacent-channel power extinction ratio” or “adjacent-wavelength power extinction ratio” or “adjacent wavelength-channel extinction ratio”, or “adjacent channel extinction ratio” for the beam from input slit SLI1 1201 to output slit SLO1 1401. It is also related to what is known to those skilled in the art as “adjacent-channel crosstalk rejection” or “adjacent-channel crosstalk extinction”. These terminologies will thus be used interchangeably below.Adjacent Channel Extinction Ratio of a Spectrometer\nLet ηace(λSM) denotes the adjacent channel extinction ratio for the spectrometer at its center operating wavelength λSM. It is called the “adjacent channel extinction ratio of a spectrometer”. This adjacent-channel extinction ratio ηace(λSM) is typically higher than about 100, especially when the size of the spectrometer is small (i.e. when the RASM factor is smaller than about 104). For many applications in photonic integrated circuits, electronic-photonic integrated circuits, and fiber-optics, ηace(λSM) higher than about 100, and RASM factor smaller than about 104 are desirable but they are largely not reachable via the methods in the prior arts.\nFurthermore, for many applications in photonic integrated circuits, electronic-photonic integrated circuits, and fiber-optics, it is desirable that at the wavelength range of interest, the efficiency of passing or detecting the input spectral power be efficient, which is given by the spectral power output efficiency factor ηeffI1-O1(λI1-O1) defined above.\nSpectral Power Output Efficiency of a Spectrometer\nLet ηeff(λSM) denote the spectral power output efficiency ηeffI1-O1(λI1-O1) when the output slit wavelength λI1-O1 1321 is at the center operating wavelength λSM of the spectrometer. In an ideal situation, ηeff(λSM) shall be equal to unity (i.e. the number “1” so that the output power is equal to the input power or it has 100% passing or detection efficiency) or at least near unity. In many applications ηace(λSM) higher than about 100 with RASM factor smaller than about 104 and ηeff(λSM) higher than about 0.5 or having a higher than 50% passing/detection efficiency is desirable. In some other applications, ηace(λSM) higher than about 100 with RASM factor smaller than about 104, and ηeff(λSM) higher than about 0.1 (i.e. having a higher than 10% passing/detection efficiency) are desirable. These capabilities are desirable for the spectrometer but they are largely not reachable via the methods in the prior arts.\nConventional Rowland Gratings and their General Limitations\nAlthough a small input slit width is required to achieve high spectral resolution for a spectrometer, a smaller input slit width WI1 1291W (see FIG. 1A), however, leads to a larger beam divergence full angle θdvdf-BI1-IP % 1141 (see FIG. 2) from the input slit defined as the beam angular span that encompasses IP %=95% of the beam energy. It can be shown that the conventional Rowland design in the prior art works reasonably well only up to (θdvdf-BI1-95%)/2<4° (for a grating in air for which the beam propagation region has a refractive index of ngr=1). More specifically, θdvdf-BI1-95% is the full divergence angle of the input beam due to spatial diffraction that we refer to as the “divergence-diffraction” (dvdf) angle. The 95% refers to the angle being defined by 95% of the beam power to be within the divergence angle. Note that the 4° is half the divergence angle of θdvdf-BI1-95%. If the grating cannot intercept or refocus a large percentage of the beam power well, then it will result in an unnecessarily large beam spot size at the output slit, and hence a worse wavelength resolution at the output slit.\nLet the input beam be a monochromatic beam with wavelength λI1-O1 1321 (see FIG. 2) and the focal spot size for this beam at the output slit SLO1 1401 be WBI1-O1-IP % 1381W (see FIG. 2). When θdvdf-BI1-95%/2>4°, the Rowland design could not give a sharp enough focus at the output slit SLO1 1401, thereby limiting the resolution of the spectrometer.\nA half-diffraction angle of θdvdf-BI1-95%/2=4° corresponds to a slit width of about 25 microns (for an input beam wavelength around λI1-O1=1000 nm and n=1). In the current art, it is typically difficult to make slit width smaller than about 25 microns, and Rowland design is adequate for most present spectrometers with slit widths larger than about 25 microns (for a grating in air for which the beam propagation region has a refractive index of ngr=1).\nSlit size (n = 3.4)θD(>95% energy)10μm 6°6μm10°2μm24°1μm46°0.5μm70°0.3μm90°\nFocusing Aberration Limitation:\nIn the case of the Rowland design, when (θdvdf-BI1-95%)/2>4 Degree(s) (DEG), serious aberration in the refocusing beam will occur to limit wavelength resolution (θdvdf-BI1-95% 1141 is the full divergence angle of the input beam due to spatial diffraction that we refer to as the “divergence-diffraction” (dvdf) angle. The 95% refers to the angle being defined by 95% of the beam power to be within the divergence angle). This is shown in FIG. 4, illustrating the ray tracing for the typical Roland-Echelle design at 4, 8, and 16 DEG diffraction. The ray tracing will allow us to see potential focusing distortion or aberration at the exit slit. In the figure, we show the focusing behavior for two sets of rays with wavelengths separated by 0.4 nm. Let the “half angle” θdiv=(θdvdf-BI1-95%)/2. From the figure, we see that their focused spots are clearly separated when θdiv=4 DEG. However, when θdiv=8 DEG, the focused spots began to smear out. FIG. 4A shows the cases of θdiv=16 DEG, and θdiv=50 DEG, for which the focused spots smear out even more. Thus, there is substantial distortion for the focusing rays when θdiv>4 DEG. Further simulations based on numerical solutions to Maxwell's wave equations using finite-difference time-domain (FDTD) method also show similar onset of focused spot size distortion at θdiv>4 DEG. In short, the current designs are close to their resolution-area (RASM) limits at θdiv of about 4 DEG and cannot be made substantially more compact without losing wavelength resolution."}
{"text": "Saw blades beveling is an advantageous capability to be included with power saw assemblies, such as table saw assemblies which include a saw assembly comprising a saw blade, extending through a planar surface (table), coupled with a motor by a mounting assembly. Beveling assemblies, provide the operator of a table saw with the ability to adjust the angular presentation of the saw blade relative to the planar surface of the table. Unfortunately, many beveling assemblies included with table saw assemblies, and the like, are limited in the capabilities they may provide to the operator.\nA typical limitation of many beveling assemblies is the enablement of beveling in a single tilt orientation. For instance, a beveling assembly may provide the operator with the ability to adjust the angular presentation of the saw blade relative to the planar surface of the table in a right tilt orientation or a left tilt orientation. Other beveling assemblies have limited beveling capabilities resulting from unwanted contact between the component features of the saw assembly or beveling assembly and other structural impediments of the table saw assembly, such as the unwanted contact between the motor, which drives the saw blade of the table saw assembly, and the underside of the planar surface of the table or structural support features of the table saw assembly. Additionally, unwanted contact may result from contact between the mounting assembly of the saw blade and the table. For instance, when the blade is beveled, the mounting assembly may contact the underside of the table because of the mounting assembly's proximate location to the table.\nUnfortunately, attempts to resolve these problems have resulted in assemblies which may require a user to accept limited capabilities. For example, unwanted contact is decreased, but only by limiting the depth of cut capabilities of the saw blade of the table saw assembly. Therefore, it would be desirable to provide a table saw assembly which enables the saw blade with bi-directional beveling capabilities (right and left tilt orientation) without reducing the capabilities of the saw blade due to factors, such as unwanted structural impedance.\nAdditional limitations may be a result of design limitations. For instance, the assembly, which couples the saw blade with the motor of the table saw, typically, statically fixes the position of the motor relative to the saw blade. Such a situation may result in decreased performance of the saw assembly during beveling due to changes in tolerances resulting from changes in the positions of these features due to beveling.\nTherefore, it would be desirable to provide a beveling assembly which increased the ability of the component features of the table saw assembly to remain near optimal tolerances during the beveling of the saw blade. Another design limitation is the typical arbor assembly is enabled as a single axle, which may be limited to coupling with the saw blade to impart angular momentum on a single end of the axle. This may result in the beveling assembly enabled with the single axle arbor assembly limiting its capabilities to a single tilt orientation, as described previously. Still further, this may result in increased time spent attempting to properly re-align the single axle arbor assembly when transitioning between beveling in a first tilt orientation to beveling in a second tilt orientation, such as from right to left. This may significantly decrease productivity accomplished through use of the table saw employing such a single axle arbor assembly. Further, due to such drastic re-positioning required by the single axle arbor assembly there may be an increase in tolerance misalignments resulting in increases in non-production time spent by the power tool. Therefore, it would be desirable to provide an arbor assembly which decreased time spent re-aligning the arbor assembly and increased the ease of employing the saw blade on the arbor assembly for beveling in both tilt orientations.\nFurther, it is common practice in the field of power tools, particularly arbor assemblies for table saws, to take steps to “true” the assembly after it has been assembled. This truing may assist in increasing the productivity and life span of the power tool by assisting in reducing wear to the component features of the power tool. Steps, such as machining the flange may true the arbor assembly and may decrease the amount of “wobble” introduced into the saw blade during operation of the table saw. Unfortunately, after the truing of the arbor assembly is complete the parts of the arbor assembly, such as the bearing for example, may wear out due to use, requiring replacement. It is known to allow for the operator to disassemble the arbor assembly in order to replace the bearing. This often requires removal of the flange in order to access the bearing. It is typically the case that when the operator re-assembles the arbor assembly the flange and possibly other component features are no longer in the trued position, as originally provided by the manufacturer. Thus, “wobble” may be introduced into the spinning blade of the table saw. Therefore, it would be further desirable to provide an arbor assembly which assists the operator in maintaining its trued position even after replacement of component features of the arbor assembly.\nTherefore, it would be further desirable to provide an arbor assembly which assists the operator in maintaining its trued position even after replacement of component features of the arbor assembly."}
{"text": "Mating fasteners that include two opposed fastener members are popular in a variety of applications. The particular type of mating fastener with which the present invention is concerned includes two opposed fastener members, each including a surface having a plurality of arranged headed stems projecting therefrom. Such fasteners are particularly useful in an environment where it is desirable to interengage the fastener using a relatively small amount of force, while requiring a larger force to separate the fastener members.\nA conventional molded polymeric fastener is shown in FIG. 1. Included are opposed fastener members 10 and 12, including bases 14 and 16, respectively. Each base includes a plurality of arranged stems 18 orthogonally projecting therefrom, and each terminating in a head 20 at the respective distal ends of stems 18. The heads have generally arcuate surfaces 22 opposite the respective bases 14 and 16, which surfaces are adapted for sliding contact with the encountered arcuate surfaces 22 of heads 20 when the fastener members are first brought together. Heads 20 also each include a generally planar latching surface 24 extending radially outwardly of the stem, which is adapted to engage a latching surface of an opposed head, as shown in FIG. 2.\nFIG. 1 shows fastener members 10 and 12 prepared for engagement, with the heads of one fastener member positioned generally between the heads of the other fastener member. To engage the fastener, fastener members 10 and 12 are urged toward each other with the respective bases generally parallel, causing heads 20 to contact and to deflect the attached stems laterally. The lateral deflection of the stems allows the heads of one fastener member to pass the heads of the opposed fastener member to engage the fastener members in the manner illustrated in FIG. 2.\nIn the past, fasteners of the type discussed above were generally formed by passing a backing layer having a multiplicity of projecting stems beneath a heat source, which heat source tends to melt the stems to form heads atop each of the stems. Such a process is known in the art, as illustrated by U.S. Pat. No. 4,290,832 (Kalleberg), entitled \"Method for Making Fasteners.\" The '832 patent discloses a multi-step process for producing a fastener member having a backing layer, a plurality of monofilament stems projecting normally from the backing layer, and a head atop each stem.\nThe process, illustrated schematically in FIG. 3, includes moving two backing layers 50 and 52 in parallel paths while feeding monofilament 54 between the layers, and pressing the monofilaments using applicators 56 and 58, so that the monofilaments extend normally between the layers. The monofilaments are severed halfway between the backing layers by blade 60, which produces dual, opposed backing layers each having a plurality of monofilaments extending from each backing. The newly severed terminal ends of the monofilaments are then heated by heat sources 62 and 64 to a temperature sufficient to form a head atop each of the stems. After the heads have solidified, the two fastener members 66 and 68 are wound onto separate rollers 70 and 72, respectively, for dispensation and use.\nWhile having their own utility, the fasteners formed by the process described above tend to exhibit inconsistently shaped heads, which can potentially affect release performance. For example, a generally spherical head having an arcuate latching surface may be unintentionally formed by such a process. An arcuate latching surface tends to disengage more easily than a planar latching surface extending radially from the stem, and could lower the force required to disengage the fastener. Alternatively, the heads could be irregularly formed due to the inconsistent application of heat, differences in production speed, and the like. Such irregularities could prompt inconsistent release characteristics, which may be undesirable under some circumstances.\nFurthermore, for certain applications it may be preferable to provide heads that are generally conical, spherical, or some other such desired shape in order to obtain certain engagement and release characteristics. For example, U.S. Pat. No. 4,454,183 (Wollman) discloses a three lobed head structure formed by the application of heat to a triangular monofilament stem. However, as discussed above, heads formed by the application of heat to a monofilament stem tend not to be consistently formed. Furthermore, the number of different head structures is generally limited by the size and shape of the stems from which they are formed.\nAnother approach to forming heads atop stems in a desired shape is disclosed in Japanese patent publication no. 1-238805 (Anzai). The method disclosed in the '805 publication involves forming individual rows of stems by injection molding. Each individual row of stems, which stems may include spherical, conical, or other shaped heads, is typically formed in a mold that is split along the longitudinal axis of the stem. After the molten material has been injected and solidified, the halves of the mold may be separated, and the individual row of stems removed.\nIn order to form a fastener from the individual rows of stems, several such rows must be aligned and connected together. The individual rows may be connected by means of adhesive, thermal bonding, high frequency bonding, and the like. Although this method of forming a fastener may represent an improvement over the prior art in terms of head formation, it involves a greater investment of time and equipment to connect the individual rows together to form the fastener, and is therefore more expensive and undesirable.\nFinally, it is often desirable to form a secondary fastening portion on one surface of the fastener, typically the back surface, in order to enable the fastener to be easily attached to a surface, body, or the like. At present, the secondary fastening portion is typically attached to the back of the fastener member by welding (either solvent, heat, or ultrasonic welding), which requires a manufacturing step separate from the formation of the fastener portion. The second manufacturing step increases the cost of the article, and is therefore undesirable.\nIt is therefore desirable to provide a method of making a unitary fastener member having rows and columns of stems having consistently identical head structures of a desired shape, to provide a method of making a unitary fastener having head structures of different desired shapes, and to provide a fastener member having an integral secondary fastening portion that is formed concurrently with the fastener member."}
{"text": "Proximity detection devices are used in a wide variety of applications for determining the relative nearness or separation of an object or person relative to another object or person. An application of such devices that has recently acquired considerable public interest involves using such devices to allow a responsible individual, such as parent or guardian, to monitor the whereabouts of another person in the custody of the responsible individual. One example of a monitoring system that is intended to equip a responsible individual with the ability to monitor the whereabouts of another person, such as a child, is described in the U.S. Pat. No. 4,598,272, to Cox, the disclosure of which is incorporated herein.\nIn accordance with the system described in the Cox patent, proximity monitoring is performed by equipping each of a parent and a child with a respective radio transmitter-radio receiver pair. The radio transmitter carried by the child is operative to broadcast or radiate a high frequency RF electromagnetic wave to which the receiver in the unit carried by the parent is tuned. Should the level of the received RF signal monitored by the parent's radio receiver drop below a prescribed threshold, indicating that the child has moved to location beyond a prescribed distance from the parent, an alarm signal is generated by the parent's device. In response to the alarm signal, the parent may then active his or her own radio transmitter, which broadcasts a radio wave that has substantial signal strength, so that it will be detected by the child's unit. When the child's unit detects the RF signal transmitted by the parent's unit, it outputs an audible alarm signal that is loud enough to be heard by the parent, thereby enabling the parent to locate the child.\nNow although the signal strength-monitoring--alarm--response radio wave communication scheme employed in the system described in the Cox patent facilitates alerting a parent/guardian of the separation of a child, patient, etc. beyond a prescribed distance, it has been found that the use of a high frequency radio wave as the signalling mechanism has a number of inherent problems that limit its performance, and thereby prevent its practical application to a wide variety of signalling environments.\nMore particularly, a fundamental shortcoming of a high frequency radio wave (e.g. one on the order of 300+ MHz) is the fact that radio waves are subject to multipath propagation, which can be especially severe in the interior of a building. Secondly, colinearity between the direction of polarization of the broadcast RF signal and the receiver antenna is required for ensuring optimum signal reception. A further problem is the effect of the dielectric distortion effect of the human body (a substantial salt water mass) on the signal, which can typically causing a fluctuation in signal amplitude on the order of 10-15 dB. Moreover, since the signal strength of a radiated electromagnetic wave is inversely proportional to the square of the distance from the emitter, the setting of a signal strength threshold to trigger an alarm yields very imprecise results, especially when considering the other effects described above."}
{"text": "The invention relates generally to mops, and more particularly to mops in which a connector is used to connect mop fibers to a handle.\nU.S. Pat. No. 1,924,978 discloses an early version of a cleaning product with a separate connector used to secure a \"dusting element\" to a handle. In that device, the connector (head H) is illustrated as being attached to the handle ferrule 5 that is anchored to the handle by a pin 6 or the like. The connector includes two deformable arms 2 that are used to hold the dusting element. The dusting element includes an internal casing or envelope 7 with a central opening 8, through which the arms can pass when squeezed together. After the arms are inserted into the casing or envelope, they can be released to hold the dusting element in place. Unfortunately, the device appears to require an internal casing or envelope, which could be unduly expensive for use with a mop.\nU.S. Pat. No. 2,299,480 shows a more conventional method for attaching mop fibers to a handle. The illustrated mop head includes a U-shaped wire loop 11 that has pointed ends 13 that can be used to pierce a cylindrical ferrule 12. The completed head can be attached to a handle by a driving a nail through a hole 14 in the ferrule.\nU.S. Pat. No. 4,135,272 discloses a more modem implementation of the idea. That discloses a plastic connector 12 having an aperture 27 and slot 29 used for securing fibers to the head. A separate strap 34 having a head 35 is first passed through the aperture 27, then under the mop fibers, then up through the slot 29. The head holds one end of the strap in the aperture 27, while projecting teeth 38 hold the other end in the slot 29. The connector includes an upper cylindrical portion 14 that permits it to be threaded onto a handle 18.\nU.S. Pat. No. 4,377,879 discloses a one-piece mop connector comprising a tongue 16 that can be used to hold mop fibers. The tongue can be locked in place by a series of grips 18 that can engage locking members 22 molded onto the sides of the connector. The connector is secured to a handle through the seating of an interior retaining projection 28 in a groove 30 on the end of the handle. This semi-permanent connection of the connector to the handle is not particularly desirable.\nIn apparent recognition of the need for improving the security of the connection of a connector to a handle, while permitting easy replacement or exchange of a mop head, U.S. Pat. No. 5,375,286 discloses a resilient bayonet-type mounting system for cleaning implements. The disclosed mop head apparently traps mop fibers between distinct upper and lower elements 9 and 10 that are secured together by ultrasonic welding.\nA simpler and more convenient connector for attaching mop fibers to a mop handle would be desirable."}
{"text": "1. Field of the Invention\nThe invention relates to rifle actions in general and bolt actions in particular.\n2. Prior Art\nIn a bolt action rifle, the barrel is typically threaded into the receiver. The barrel includes the chamber, which is that portion of the barrel that holds the round immediately prior to firing. The rough dimensions of the chamber are formed prior to the barrel being threaded into the receiver. However, the final dimensions of the chamber are formed by hand in order to ensure that the headspace is proper.\nHeadspace is the distance from the face of the bolt to whatever surface in the chamber stops the round from advancing into the barrel. In a rimless cartridge, such as the .308, the headspace is the distance from the bolt face to the shoulder in the chamber. Stated more broadly, in a rimless cartridge, the headspace is a function of the distance from the face of the bolt to the chamber and of the chamber dimensions. Headspace in a rimless cartridge is illustrated in FIG. 1 as dimension X.\nThe acceptable tolerance for headspace in rifles is relatively narrow. In a .308, the acceptable headspace range is 1.6300 to 1.6340 inches. Thus, only 4 thousandths of an inch separate too little headspace from too much.\nVery bad things can happen if the headspace is not within specification for the rifle. When the round is fired, the chamber is supposed to hold the casing in place so that the expanding (i.e., exploding) gases can be used to drive the bullet down the barrel. However, if there is too much headspace, those same gases can drive the casing back toward the bolt. This can result in failure of the casing, damage to the bolt, and in some cases, even catastrophic failure of the rifle. Excessive headspace can also result in the round advancing into the chamber so that the firing pin in the bolt cannot adequately impact the primer, with the result that the round fails to go off. Similarly, excessive headspace can allow the spent shell casing to slide forward into the chamber so that it cannot be extracted, which can preclude the rifle from being reloaded—an annoyance in the best of circumstances and a potentially deadly occurrence to the shooter in certain conditions, such as self-defense applications and dangerous game hunting.\nInsufficient headspace can be a serious problem as well. If the bolt is closed on a round in a rifle with insufficient headspace, the bolt may be exerting tension on the casing in the chamber. This can lead to casing failure, resulting in the emission of the hot expanding powder gases into the chamber. This can result in catastrophic failure of the rifle and/or emission of hot gases directly into the face of the shooter. It can also result in a portion of a failed casing not extracting from the chamber.\nBecause of the close tolerances required of headspace and the significant risks associated with the operation of a rifle whose headspace is out of specification, the chambering of a rifle is typically finished by hand. First, the rifle will be rough chambered. Then, the gunsmith installing the barrel must manually ream the chamber, very slowly. He turns the reamer, measures the headspace, and then turns the reamer some more, until the headspace measurement is correct. This is a very time consuming process, and it is a substantial factor in the cost of rifle manufacturing as well as after market barrel replacement.\nIn view of the foregoing, a bolt action rifle meeting the following objectives is desired."}
{"text": "Medical devices may be used during internal operations on a patient. During these operations, bodily fluids, such as blood, and/or tissue may come into contact with the medical device. Once blood and/or tissue come into contact with the medical device, the device may lose its sterility and may be contaminated. In general, medical devices may be constructed with little consideration for the potential of disassembly since the medical devices may be designed to be disposed of once contaminated after a single use. In some instances where portions may be made to be removable, they may generally consist of disposable parts for the device that detach to be thrown away before a new disposable part is used for a new patient. For some medical instruments lacking disposable parts, they may generally either be only used once or may need to be entirely resterilized before reuse.\nWith the advancement of the electronics industry, medical devices may be adapted to contain most, if not all, of the required components within the medical device. More specifically, some medical devices may be adapted to use an internal or attachable power source instead of requiring the device to be plugged into an external source by a cable. Merely exemplary devices that may be adapted to include a portable power source are disclosed in U.S. Pat. No. 6,500,176 entitled “Electrosurgical Systems and Techniques for Sealing Tissue,” issued Dec. 31, 2002, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,416,101 entitled “Motor-Driven Surgical Cutting and Fastening Instrument with Loading Force Feedback,” issued Aug. 26, 2008, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,738,971 entitled “Post-Sterilization Programming of Surgical Instruments,” issued Jun. 15, 2010, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2006/0079874 entitled “Tissue Pad for Use with an Ultrasonic Surgical Instrument,” published Apr. 13, 2006, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2007/0191713 entitled “Ultrasonic Device for Cutting and Coagulating,” published Aug. 16, 2007, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2007/0282333 entitled “Ultrasonic Waveguide and Blade,” published Dec. 6, 2007, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2008/0200940 entitled “Ultrasonic Device for Cutting and Coagulating,” published Aug. 21, 2008, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2009/0209990 entitled “Motorized Surgical Cutting and Fastening Instrument Having Handle Based Power Source,” published Aug. 20, 2009, issued as U.S. Pat. No. 8,657,174 on Feb. 25, 2014, the disclosure of which is incorporated by reference herein; and U.S. Pub. No. 2010/0069940 entitled “Ultrasonic Device for Fingertip Control,” published Mar. 18, 2010, the disclosure of which is incorporated by reference herein. Similarly, various ways in which medical devices may be adapted to include a portable power source are disclosed in U.S. Provisional Application Ser. No. 61/410,603, filed Nov. 5, 2010, entitled “Energy-Based Surgical Instruments,” the disclosure of which is incorporated by reference herein.\nElectrically powered medical devices such as those referred to herein may require the use of high value or environmentally restricted disposable components to operate. The ability to reuse or reprocess these components over multiple uses may increase the value of the initial purchase by possibly spreading the cost of those components over the multiple uses. It may also be desirable in some settings to allow portable electronic components such as batteries to be recharged and/or be otherwise reprocessed between uses in sterile medical devices, without such electronic components having to be resterilized between uses in sterile medical devices, and without such devices contaminating sterile medical devices during re-use of the non-sterile electronic components. In addition or in the alternative, reclamation and reuse of the components may avoid or mitigate any environmental issues that may otherwise be associated with the disposal of the components after a single use. One potential approach may include reusing clean electrical components within multiple medical devices.\nWhile several systems and methods have been made and used for component acceptance and release features for medical devices, it is believed that no one prior to the inventors has made or used the invention described in the appended claims.\nThe drawings are not intended to be limiting in any way, and it is contemplated that various embodiments of the technology may be carried out in a variety of other ways, including those not necessarily depicted in the drawings. The accompanying drawings incorporated in and forming a part of the specification illustrate several aspects of the present technology, and together with the description serve to explain the principles of the technology; it being understood, however, that this technology is not limited to the precise arrangements shown."}
{"text": "This invention relates to an arrangement for coupling undeveloped films to a leader.\nWhen developing undeveloped films (hereinafter simply referred to as \"films\") in an automatic film developing machine of a type that utilizes leaders to guide films, the front ends of films F are connected to a leader L after pulling them out of patrones P as shown in FIG. 27. The films F are sent to the film developing unit, led by the leader.\nThe leader L is formed of a flexible synthetic resin sheet. Along its longitudinal center are formed a plurality of holes a arranged at equal intervals and adapted to engage a feed sprocket provided in the film developing unit.\nSuch an automatic film developing machine is ordinarily provided with a film guide. If films are detached from the film guide while being fed, their surfaces may be damaged due to contact with the film guide. Thus, it is necessary to couple films reliably to the leader with high positional accuracy.\nIf a film F should come off the leader L while developing, it will be stuck in the treating solution. It is extremely troublesome to take out such a stuck film. If taken out carelessly, the film F may be exposed to light. Thus, films F have to be rigidly coupled to leaders L.\nFor this purpose, films F were heretofore connected to a leader F by means of splicing tapes T as shown in FIG. 27.\nIn order to connect films to a leader using splicing tapes T with high positional accuracy, the leader has to be placed on a special-purpose workbench. It is necessary to keep oil and dust from the adhesive surface of the splicing tapes T. Thus, a taping job tends to be very troublesome. In order to prevent films F from coming off the leader while being fed, it is necessary to use splicing tapes with high bond strength. It is difficult to peel off such tapes after developing films F. Such tapes are not recyclable and thus are uneconomical."}
{"text": "1. Field of the Invention\nThis invention lies in the field of sulfur removal and recovery, and particularly in the treatment of sour gases and other gases in the chemical process industry that contain hydrogen sulfide.\n2. Description of the Prior Art\nMuch of the natural gas produced in the United States has a hydrogen sulfide content exceeding 4 parts per million and is therefore classified as xe2x80x9csour.xe2x80x9d Since hydrogen sulfide is an environmental hazard, sour natural gas is considered unacceptable for transport or use. Hydrogen sulfide levels are also a problem in the fabrication of fuels derived from petroleum, coal and lignite, whose sulfur content is converted to hydrogen sulfide during the conversion of these materials to gasoline, jet fuels, refinery gas, coal gas, blue-water gas and the like. In addition to the environmental hazard, hydrogen sulfide represents a loss of the sulfur value of the raw material, which if recovered as elemental sulfur would be of significant use to the chemical industry.\nThe traditional method for converting hydrogen sulfide in natural gases and in gaseous plant effluents is the Claus process, in which part of the hydrogen sulfide is burned in air to form sulfur dioxide and water:\n2H2S+3O2xe2x86x922SO2+2H2Oxe2x80x83xe2x80x83(A)\nand the sulfur dioxide thus produced is reacted with further hydrogen sulfide to form sulfur and additional water:                                           2            ⁢                          xe2x80x83                        ⁢                          H              2                        ⁢            S                    +                      SO            2                          ->                                            3              x                        ⁢                          S              x                                +                      2            ⁢                          xe2x80x83                        ⁢                          H              2                        ⁢            O                                              (        B        )            \nThe symbol xe2x80x9cxxe2x80x9d in Reaction B is used to denote that the elemental sulfur exists in a mixture of molecular species varying in the number of sulfur atoms per molecule.\nThe furnace (Reaction A) in the Claus process is operated with a fuel-rich mixture, converting only one-third of the H2S to SO2. The fuel-rich atmosphere results in the partial conversion of hydrocarbons that are present in the H2S feed to such compounds as COS and CS2, which lessen the yield of elemental sulfur and are themselves hazardous. The fuel-rich atmosphere also promotes the breakdown of aromatics to soot. For high sulfur recovery, precise control of the overall stoichiometry is needed, and this is made especially difficult when considerable amounts of CO2 and other inerts are present.\nPart of Reaction B occurs in the furnace and the rest is conducted in a heterogeneous system in which the reaction mixture is gas-phase and contacts a solid activated alumina catalyst of a sort well known to those skilled in the art of the Claus process. With continued use, the alumina catalyst fouls and becomes otherwise deactivated over time. This requires plant shutdown, loss of process time, and the cost of regeneration or replacement of the catalyst, together with the associated labor costs.\nA further disadvantage of Reaction B is that it is equilibrium-limited at temperatures above the dewpoint of sulfur, and despite being performed in two to four stages, the reaction leaves 2% to 5% of the H2S and SO2 unreacted. Each stage requires a separate condenser to remove the elemental sulfur, and these condensers require a large heat-exchange area and reheating of the gas leaving each but the last condenser. Furthermore, the steam generated by each condenser is low in pressure, limiting its usefulness. Additional costs are entailed in treating the tail gas in which the sulfur content must be reduced by ten to twenty times.\nIt has now been discovered that virtually complete conversion of hydrogen sulfide in natural gas or other gas mixtures to elemental sulfur and water can be achieved with the use of a single-stage reaction between hydrogen sulfide and sulfur dioxide, in a manner producing no reaction products other than elemental sulfur and water. The reaction\n2H2S+SO2xe2x86x923S+2H2Oxe2x80x83xe2x80x83(I)\nis conducted with excess H2S in the liquid phase in the presence of a homogeneous liquid-phase Claus catalyst, at a temperature above the melting point of sulfur but low enough to keep the reaction in the liquid phase, and upstream of the furnace where unreacted H2S is combusted to produce the SO2 that is consumed in reaction (I). The H2S that is combusted in the furnace is the excess H2S that passes unreacted through the liquid-phase reaction, optionally supplemented by H2S from an H2S-containing stream that bypasses the reaction (I) reactor. The SO2 in the furnace combustion gas is recycled to the reaction (I) reactor either as a gas or dissolved in a solvent, and in any case serves as the entire SO2 feed to the reaction.\nIn basic terms, the invention as shown in FIG. 1 proceeds as follows:\n(a) In a first stage 1, an H2S-containing mixture is passed through a continuous-flow catalytic reactor where the mixture contacts SO2 in accordance with the reaction\n2H2S+SO2xe2x86x923S+2H2Oxe2x80x83xe2x80x83(I)\nusing approximately 10% to 50% excess H2S. The H2S enters the reactor either as a gas or dissolved in an organic solvent; in most cases the H2S will enter as a gas. The SO2 likewise enters the reactor either as a gas or dissolved in an organic solvent. Regardless of the phases of the streams entering the reactor, both reactants dissolve in an organic liquid solvent flowing through the reactor, the solvent entering either with one of the two incoming reactant streams with the H2S or the SO2 dissolved in the solvent, or as a circulating stream recycled from the reactor exit. The reaction causes a major fraction of the SO2, preferably substantially all of it, to react. The term xe2x80x9ca major fractionxe2x80x9d is used herein to indicate at least half, and preferably 80-90% or more. The organic liquid solvent also contains a dissolved catalyst that promotes reaction (I). The reaction produces elemental sulfur, which is recovered from the product mixture by phase separation. The reaction is conducted at a temperature above the melting point of sulfur and below the boiling point of the solvent, preferably below the temperature at which sulfur polymerizes.\n(b) In a second stage, 2, H2S, including H2S that passed unreacted through the first stage, is combusted with oxygen according to the reaction\nxe2x80x832H2S+3O2xe2x86x922SO2+2H2Oxe2x80x83xe2x80x83(II)\nto convert the H2S to SO2. Hydrocarbons that may accompany the H2S are combusted to CO2 and H2O whereas organic sulfur compounds that may be present additionally yield SO2.\n(c) In a third stage, 3, the SO2 produced in the second stage is recovered by absorption and returned to the reactor (the first stage). The SO2 in the third stage may be recovered in the solvent used in the first stage, 1, with the first stage catalyst dissolved in the solvent. The solution containing both the SO2 and the catalyst can then be recycled in its entirety to the continuous-flow reactor (the first stage) as the SO2 feed to the reactor. Alternatively, the solvent used to recover the SO2 in the third stage, 3, may be kept separate from the solvent in the reactor of the first stage, 1, with the SO2 being stripped from it and sent to the first stage, 1, as a gas. The third stage absorption leaves a tail gas that is substantially free of H2S and SO2.\nVarious additional process stages upstream, intermediate and downstream of these three stages are included in any of various arrangements in preferred embodiments of the invention to enhance the flow and transfer of streams, to separate phases and control concentrations and flow rates, to separate the water formed in Reactions I and II, and to control other process parameters such as temperature and pressure. These and. other characteristics, features and advantages of the invention will be better understood from the description that follows."}
{"text": "1. Technical Field\nThe present disclosure generally relates to optical transceivers, and more particularly to a packaging method of an optical transceiver.\n2. Description of Related Art\nOptical communications has become increasingly popular and is widely used in information transmission. Generally, optical transceivers are used in optical communication devices to perform a conversion between electrical signals and optical signals. An optical transceiver comprises a plurality of laser components and a plurality of lenses. Definition of position between the laser components and the lenses affects performance of transmitting information of the optical transceiver. Therefore, how to make an exact definition of position between the laser components and the lenses is an important problem in study.\nTherefore, a need exists in the industry to overcome the described limitations and reduce the size of the electronic device."}
{"text": "1. Technical Field\nThis invention is directed generally to devices and methods for facilitating bone growth, and, in particular, to bone implantable devices and implantation methodologies that augment beneficial bone growth or repair while limiting bone growth in undesirable directions.\n2. Background\nIn orthopedic and neurological surgical procedures it is often important to facilitate the growth or fusion of bony structures. This may entail growth “bone-to-bone” or, depending on the nature of the procedure, bone to device.\nChronic back problems, for example, cause pain and disability for a large segment of the population. In many cases, such problems are attributable to relative movement between vertebrae in the spine. Spinal surgery includes procedures to stabilize adjacent vertebrae. Common stabilization methods often involve fusing adjacent vertebrae together.\nFusion techniques include removing disc material that separates the vertebrae and impacting bone into the disc area. The impacted bone fuses with the bone material of the two adjacent vertebrae to thereby fuse the vertebrae together.\nIn a further advance in the art, spinal implants have been developed to increase the probability of a successful fusion. Such devices generally comprise a hollow cage into which bone growth inducing substances, such as bone chips or bone slurry, may be placed. The cage is inserted, either by anterior or posterior approach, into the intervertebral disc space. The cage wall has holes extending radially therethrough, typically throughout the entire cage surface. Bone growth extends into the device through the radial apertures, facilitating arthrodesis between the adjacent vertebral bone structures and allowing for the decompression of neural elements.\nWith the continued development of techniques for achieving spinal fusion through the use of spine fusion cages, new materials have been developed to augment the fusion process. Traditionally, the patient's own bone, or cadaver bone, was used in the cage to promote bony fusion. More recently, powerful new biologic materials have been discovered that greatly accelerate the fusion process, in some cases eliminating the need for donor bone.\nHowever, with the utilization of the newer biologic materials there has arisen a significant problem. When bone growth inducing agents, such as bone morphogenic proteins (“BMP”), are used in cages of existing design there is risk of inducing the overgrowth of bone around and into sensitive neural tissues. This is especially the case when a posterior approach is utilized to implant a spinal fusion cage, as bony overgrowth toward the central canal or neural foramen may impinge on spinal nerve roots causing neurological damage. A recent study on posterior lumbar interbody fusion procedures using rhBMP-2 reported that 58% of patients experienced greater than expected bone formation dorsal to the fusion cage. In 30% of the cases, the bony overgrowth compromised the central canal, the neural foramen, or both. This study is confirmatory to observations first made by the present inventor in early 1999.\nTypically, the bone growth agent is in liquid form and is applied to an absorbent carrier material, such as a piece of bovine collagen. The doped carrier material is placed with forceps into the interbody space, usually into an open end of the fusion cage after the cage has been implanted, but sometimes prior to cage implantation. During placement carrier material may inadvertently wipe across body areas, including internal bony structures, where bone growth is not desired. In addition, as the carrier material is pressed into place agent may squeeze out and flow into adjacent areas. Exacerbating the problem, current protocols do not encourage the use of suction, irrigation and hemostatic agents when bone growth agent is utilized. Conventional cage design also allows for the leakage of agent into undesirable areas after implantation through ill-placed apertures in the cage body, in the cage end caps, or otherwise. Because of the powerful stimulatory effects of bone growth agents, uncontrolled application of these substances may lead to serious complications, including severe inflammation, debilitating neural impingement, and other potential complications.\nThus, there is a need to better control the bone growth process when using a bone implantable device, especially in circumstances where powerful bone growth inducing agents are used in conjunction therewith.\nIn satisfying this need, there is also an opportunity to extend the application of bone growth agent based bony fusion to all types of bone implantable devices to better achieve union of bone-to-bone or bone-to-device, as the case may be.\nFurthermore, it is desirable to eliminate the use prior art rod and screw structure that tends to interfere with a patient's musculature and tends to kill nerves and destroy segmented branch nerve extensor muscles in patients."}
{"text": "A Web application may include HyperText Markup Language (HTML) pages, JavaServer Pages (JSPs), servlets, resources, source files, etc. A Web archive (WAR) file is a packaged Web application that may be used to deploy a Web application into a Web server."}
{"text": "1. Field of the Invention\nThe present invention relates to an X-ray computed tomography (CT) apparatus that collects projection data by exposing X-rays to a subject and generates an image based on the projection data, and relates to a method of controlling the X-ray CT apparatus\n2. Description of the Related Art\nThe X-ray CT apparatus includes an X-ray tube and an X-ray detector, which are arranged opposite to one another by sandwiching a subject. The X-ray CT apparatus rotates the X-ray tube and the X-ray detector about the subject while exposing X-rays from the X-ray tube and detecting, in the X-ray detector, X-rays transmitted through the subject. The data detected by the X-ray detector is collected as projection data in a data acquisition system (DAS), and the collected projection data is used to generate an image of the subject. For example, the projection data is collected by rotating the X-ray tube and the X-ray detector by 360° or (180°+fan angle) and the image is generated based on the projection data over 360° or (180°+fan angle).\nA set of detection data detected at an angle relative to the subject by the X-ray detector having a plurality of detection elements is referred to as a “view”. A single rotation of the X-ray tube and the X-ray detector about the subject to collect the projection data of a plurality of views necessary to reconstruct an image is referred to as a “scan”. For example, when the projection data for 1 view is collected for each 1°, the projection data for 360 views is obtained for 1 scan and the image is reconstructed using the projection data of these 360 views.\nThe collection of the projection data is continuously performed by the data acquisition system (DAS) while X-rays may be intermittently exposed from the X-ray tube. When the image is reconstructed using the projection data collected by such imaging, it is necessary to reconstruct the image using the projection data collected over a period during which X-rays have actually been exposed. For example, in cardiac testing, imaging may be performed by switching the X-ray exposure on and off according to the subject's heartbeat (prospective gating scan) (e.g., Japanese Unexamined Patent Application Publication No. 2009-28065). In this case, it is necessary to reconstruct the image using the projection data collected by the data acquisition system (DAS) over a period during which X-rays have actually been exposed.\nIn the prior art, an X-ray detector (Ref detector) that is to be a reference is provided to the X-ray CT apparatus and a detection value (count value) detected by the Ref detector is added to the projection data for each view. Then, based on the count value detected by the Ref detector and added to each projection data for 1 view, it is determined whether X-rays have been exposed and it is thereby determined whether the projection data has been collected during a period in which X-rays have been exposed. In addition, when the exposure amount of the subject during the actual imaging is obtained, the dose is calculated by counting the views for which X-rays have been exposed.\nIn addition, in a scan mode in which X-rays are exposed while modulating the X-ray output (mA), the level of X-ray output (mA) can be judged by a detection value detected by the Ref detector but the actual X-ray output value (mA value) cannot be judged. Therefore, the dose is calculated using the X-ray output value that has been set by the X-ray CT apparatus during imaging planning.\nIn the prior art, by setting a threshold relative to the detection value (count value) of the Ref detector, it is determined whether the projection data for each view has been collected during a period in which X-rays have been exposed. However, the count value of the Ref detector is varied according to each Ref detector. In addition, the count value of the Ref detector is varied according to the conditions for X-ray exposures. Therefore, it is difficult to set a threshold suitable for determining whether X-rays have been exposed or not. For example, it is necessary to take into consideration the individual differences of the Ref detectors and the variations in the individual differences of the X-ray tube. In addition, the count value of the Ref detector is varied according to the tube voltage. Therefore, it is difficult to set a suitable threshold. In addition, by taking into consideration the offset value of the Ref detector and the data acquisition system (DAS), it is necessary to decide the threshold for determining the ON and OFF states of the X-ray exposure. Moreover, it is possible to make an error in determining the ON and OFF states of the X-ray exposure if the X-ray tube is unstable and may, for example, instantaneously generate a discharge or undergo current fluctuations.\nIf there is an error in determining the ON and OFF states of the X-ray exposure, the number of views necessary for reconstruction may not be obtained or the image may be reconstructed using the projection data of views that have been collected when the X-rays were in the OFF state.\nFurthermore, using the detection value (count value) of the Ref detector makes it possible to determine only the ON and OFF states of the X-ray exposure. In a scan mode in which the X-ray output (mA) is modulated, the actual X-ray output value cannot be judged. Therefore, it is difficult to accurately calculate the actual exposure dose."}
{"text": "Field of the Invention\nEmbodiments of the present invention relate generally to computer processing and, more specifically, to automated techniques for retrofitting devices.\nDescription of the Related Art\nDecreases in prices and increases in the availability of electronics have led to a proliferation of “smart” devices. Smart devices implement more powerful and flexible functionality and/or connectivity than “legacy” devices. For example, a smart thermostat can be configured to control a heating system based on sensor inputs or via a smartphone. By contrast, legacy thermostats are typically configured to control a heating system based on only a single temperature setting that is entered manual via control buttons situated on the thermostat. Although many users would benefit from the additional functionality and connectivity that come with smart devices, users oftentimes do not replace functioning legacy devices with smart devices for various practical or emotional reasons, such as the cost of purchasing the smart devices, the time and effort to replace the legacy devices, and/or emotional attachments to the legacy devices, to name a few. Instead, in many instances, users attempt to retrofit their legacy devices to enable the legacy device to implement additional or different functionality or connectivity.\nTo retrofit a given legacy device, a user typically designs a proxy interface that provides the desired functionality and connectivity when attached to the legacy device. For example, the proxy interface could include actuators that are configured to manipulate legacy controls based on inputs that are received from a smartphone. One drawback of this approach, however, is that designing a proxy interface is typically a complex, multi-step process. First, the user has to define the high-level behavior for the proxy interface. The user then has to identify and purchase multiple component instances, such as sensors, actuators, microcontrollers, etc. Subsequently, the user has to design and manufacture an enclosure that houses the component instances and, when attached to the legacy device, enables the component instances to manipulate the legacy controls to execute the desired high-level behavior. Lastly, the user has to assemble the component instances into a circuit. In addition, if the proxy interface includes a programmable instance, like a microcontroller, then the user has to write, compile and upload the firmware used to control the operation of the circuit through the programmable instance.\nAs the foregoing illustrates, designing a proxy interface requires significant knowledge across a range of technical areas, such as enclosure design and fabrication, circuit design, and programming. A lack of knowledge in one or more of these areas can discourage a user from attempting to retrofit a legacy device. Further, because the design and assembly process involves many different, manual operations, designing proxy interfaces can be tedious and time consuming, irrespective of the technical expertise of the user.\nAs the foregoing illustrates, what is needed in the art are more effective techniques for retrofitting legacy devices with new features and functionalities."}
{"text": "Venous thromboembolism (VTE), including deep vein thrombosis (DVT) and pulmonary embolism (PE), is a common cardiovascular disorder that affects up to 2 million patients in the U.S. each year. VTE is a significant risk in surgical patient populations where preoperative, operative, and postoperative immobilization may lead to blood stasis. VTE, when poorly treated, may lead to the development of post-thrombotic syndrome (PTS), with symptoms including chronic leg pain, swelling, and ulcers. As a consequence of PTS, many individuals struggle with a reduced quality of life.\nCurrent treatment options for VTE include anticoagulation, catheter-direct thrombolysis (CDT) and pharmacomechanical thrombectomy. However, the efficacy of these methods is disputed. One deficiency of these methods is the use of thrombolytic agents, which may present a risk of internal bleeding, including the potential for an intracranial bleeding event. In particular, such agents are effective because they cause blood thinning and prevent coagulation. However, their use can cause bleeding complications. The use of methods involving thrombolytic agents may also be ineffective to treat against blood clots in larger veins, such as veins located in a patient's leg. Many thrombolytic agents are also costly, and logistically complicated to use.\nAs an alternative, mechanical compression devices that compress or squeeze portions of a patient's body to force blood flow have been used. But such devices include their own drawbacks. Oftentimes, these kinds of devices are bulky and may impose a considerable burden on hospital staff. Patients are also less inclined to ambulate while using these devices because they typically are burdensome to remove and reapply. Many patients also dislike using mechanical compression systems because such devices are typically uncomfortable. In addition, these kinds of compression devices also may not be effective at treating thrombosis in larger veins.\nAccordingly, there is a need for a device that is designed to effectively treat thrombosis in large veins without using thrombolytic agents."}
{"text": "1. Field of the Invention\nThe invention relates to the area of electronic design automation (EDA), and more particularly to optical proximity correction of sub-micron mask designs.\n2. Description of Background Art\nA fabrication mask or reticle is used when fabricating a semiconductor device. Such a mask has a light-transparent portion (e.g., glass) and a light-shielding portion (e.g., chromium) that define a circuit pattern to be exposed on a wafer. Such patterns may define diffusion regions or field oxidation regions provided on a substrate. They may also define gate electrode patterns at a polysilicon level or a metal line patterns at any metallization layer of a chip. The mask and reticle must have a precisely defined circuit pattern. Photolithography is a process used for patterning semiconductor wafers during the manufacture of integrated circuits, such as application specific integrated circuits (ASICs).\nThe reticle is placed between a radiation source producing radiation of a pre-selected wavelength and a focusing lens which may form part of a xe2x80x9cstepperxe2x80x9d apparatus. Placed beneath the stepper is a resist covered silicon wafer. When the radiation from the radiation source is directed onto the reticle, that fraction of the radiation passing through the glass projects onto the resist covered silicon wafer. In this manner, an image of the reticle is transferred to the resist. For further information on IC fabrication and resist development methods, reference may be made to a book entitled Integrated Circuit Fabrication Technology by David J. Elliott, McGraw Hill, 1989.\nLight passing by the edge of a reticle pattern feature (e.g., the boundary between a chromium coated region and a transparent region) will be diffracted, so that rather than producing a very sharp image of the feature edge, some radiation diffracts beyond the intended image boundary and into the dark regions. Hence feature shapes and sizes deviate somewhat from the intended design. For ultraviolet radiation in common use today, the intensity of the diffracted radiation drops off quickly over a fraction of a micron, so the affect does not prove particularly problematic when devices have dimensions on the order of 1 micrometer. However, as device dimensions have shrunk to the submicron domain, diffraction effects can no longer be ignored.\nThe diffraction errors can be compensated for by increasing the thickness of various critical features on the pattern. For example, increasing the width of a line on the pattern will reduce the diffraction effects. Unfortunately, this defeats the purposes of using small critical dimension features; greater logic density and improved speed.\nFIG. 1A shows a hypothetical reticle 100 corresponding to an IC layout pattern. For simplicity, the IC pattern consists of three rectangular design features. A clear reticle glass 110 allows radiation to project onto a resist covered silicon wafer. Three rectangular chromium regions 102, 104 and 106 on reticle glass 110 block radiation to generate an image corresponding to intended IC design features.\nFIG. 1B illustrates how diffraction and scattering affect an illumination pattern produced by radiation passing through reticle 100 and onto a section of silicon substrate 120. As shown, the illumination pattern contains an illuminated region 128 and three dark regions 122, 124, and 126 corresponding to chromium regions 102, 104, and 106 on reticle 100. The illuminated pattern exhibits considerable distortion, with dark regions 122, 124, and 126 having their corners rounded and their feature Widths reduced. Other distortions commonly encountered in photolithography (and not illustrated here) include fusion of dense features and shifting of line segment positions. Unfortunately, any distorted illumination pattern propagates to a developed resist pattern and ultimately to IC features such as polysilicon gate regions, vias in dielectrics, etc. As a result, the IC yield is degraded or the reticle design becomes unusable.\nTo remedy this problem, a reticle correction technique known as optical proximity correction (xe2x80x9cOPCxe2x80x9d) has been developed. Optical proximity correction involves adding dark regions to and/or subtracting dark regions from a reticle design at locations chosen to overcome the distorting effects of diffraction and scattering. Typically, OPC is performed on a digital representation of a desired IC pattern. First, the digital pattern is evaluated with software to identify regions where optical distortion will result. Then the optical proximity correction is applied to compensate for the distortion. The resulting pattern is ultimately transferred to the reticle glass. OPC is described generally at the end of this document.\nFIG. 1C illustrates how optical proximity correction may be employed to modify the reticle design shown in FIG. 1A and thereby better provide the desired illumination pattern. As shown, a corrected reticle 140 includes three base rectangular featuresxe2x80x94142, 144, and 146xe2x80x94outlined in chromium on a glass plate 150. Various xe2x80x9ccorrectionsxe2x80x9d have been added to these base features. Some correction takes the form of xe2x80x9cserifsxe2x80x9d 148a-148f and 149a-149f. Serifs are small appendage-type addition or subtraction regions typically made at corner regions on reticle designs. In the example shown in FIG. 1C, the serifs are square chromium extensions protruding beyond the corners of base rectangles 142, 144, and 146. These features have the intended effect of xe2x80x9csharpeningxe2x80x9d the corners of the illumination pattern on the wafer surface. In addition to serifs, the reticle 140 includes segments 151a-151d to compensate for feature thinning known to result from optical distortion.\nFIG. 1D shows a hypothetical xe2x80x9ccorrectedxe2x80x9d illumination pattern 160 produced on a wafer surface 160 by radiation passing through the reticle 140. As shown, the illuminated region includes a light region 168 surrounding a set of dark regions 162, 164 and 166 which rather faithfully represent the intended pattern shown in FIG. 1A. Note that the illumination pattern shown in FIG. 1B of an uncorrected reticle has been greatly improved by use of an optical proximity corrected reticle.\nOPC, as now practiced, involves modifying a digital representation of a reticle design such as that shown in FIG. 1A. The modification is performed by a computer such as workstation having appropriate software for performing OPC. Points separated by less than the critical dimension on the design are evaluated in sequence and corrected as necessary. Evaluation of each point requires analysis of surrounding features in two-dimensions to determine whether problematic diffraction effects are likely. If so, an appropriate correction (serif or segment removal, for example) is performed.\nA problem with using OPC when performing a full mask design correction is that a substantial commitment must be made in terms of time and computing power in order to optically correct the integrated circuit design, such as an ASIC design. For example, a moderately complex integrated circuit design may require at least a few days to correct with OPC even when the OPC algorithm runs on the fastest modern workstations. Often an ASIC will be attractive to a customer only if it can be designed in a relatively short period of time. If the time committment is too great, then other integrated circuits such as programmable logic devices may look more appealing. Further, an ASIC designer may need to perform a full-mask design OPC on tens of thousands of ASIC designs annually. Accordingly, a significant time and computing commitment must be invested in order to perform OPC on all of a designer\"\"s ASIC designs.\nThe computational expense of OPC can be understood by recognizing as pointed out that the correction often involves adding multiple small serifs to corners of design features and removing from, adding to, or displacing lateral sections of lines. First, these many small modifications greatly increase a pattern\"\"s complexity. Second, the modifications are made by evaluating an initial pattern with very fine granularityxe2x80x94typically evaluating potential correction points separated by no more than about 0.02 micrometers (using a 0.25 micron critical dimension technology). Note that a typical reticle design may include about 50-100 million xe2x80x9crectanglesxe2x80x9d of average size 0.5 by 0.5 micrometers. Finally, each correction is made by evaluating the surrounding pattern features in two dimensions. For example, a decision as to whether a point under consideration should be corrected may be made only after a five by five grid of surrounding points is first evaluated.\nWhat is needed is a system and method for efficiently and accurately performing OPC on integrated circuit designs.\nThe invention provides a system and method for performing optical proximity correction on an integrated circuit mask design by initially performing optical proximity correction on a library of cells that are used to create the integrated circuit. The pre-corrected cells are imported onto a mask design. All cells may be placed a minimum distance apart (e.g., about 3 to 5 times the wavelength of the photolithography light) to ensure that minimal degradation from proximity effects will occur between elements fully integrated in different cells. Alternatively, the cells may be packed more closely so long as the perimeter or interfacial regions of the cells are subjected to OPC: preferably one-dimensional OPC. Because this OPC step is limited to the perimeter or interfacial regions and is performed in one-dimension, the total computational investment in correction for any given design is minimal.\nOne aspect of the invention involves a method of designing an integrated circuit that may be characterized as including the following: (a) selecting optical proximity corrected cells from one or more libraries of such cells; (b) placing one or more instances of cells selected from the one or more libraries adjacent to one another in an integrated circuit design; and (c) performing optical proximity correction on an integrated circuit mask layout at an interface between the instances of the cells on the integrated circuit design. The xe2x80x9cinstancesxe2x80x9d of the cells refer to individual copies of the cells as they exist as residents of an integrated circuitxe2x80x94as opposed to as generic versions generally available in a library.\nThe instances of the selected cells should be closely placed on the design, preferably no further from one another than about 3 to 5 times the wavelength of radiation to be used in performing photolithography during fabrication of the integrated circuit. After they are so placed, routing lines should be provided between them. It is these routing lines, outside the domain of cell, where optical proximity correction is typically performed.\nWhile one-dimensional optical proximity correction on the interfaces between the cell instances is preferred, the invention may also be practiced with two-dimensional OPC.\nAnother aspect of the invention involves computer based systems that may be characterized as including the following items: (a) a storage device; (b) a cell library, located in the storage device, and comprising a plurality of cells, each containing a predefined integrated circuit functional component; (c) a place and route tool for placing instances of specified cells, selected from the library, on an integrated circuit mask layout; and (d) an optical proximity correction module for performing optical proximity correction on interfacial regions of the specified cells as provided on the integrated circuit mask layout.\nPreferably, the place and route tool is configured such that cells in the mask design are separated by a distance of no greater than about 3 to 5 times the wavelength of radiation to be shown through a reticle containing the integrated circuit mask layout during photolithography. Preferably, the optical proximity correction module is configured to perform one-dimensional optical proximity correction on the interfacial regions containing features not appearing in the specified cells.\nThis invention also relates to reticles containing corrected layouts as produced by the systems and methods of this invention. Such reticles are intended for use in photolithography during which the corrected design layout on the reticle is used to produce an image on a wafer under fabrication. The invention further relates to integrated circuits fabricated with corrected designs produced by the systems and methods of this invention.\nThe invention may be further understood with reference to the following detailed description and associated drawings."}
{"text": "The present invention relates to a thiol-cured epoxy composition that is particularly useful as a sealer for a joint between two substrates.\nJoints often are coated with a sealer or sealant to protect the joint against the penetration of water, oil, salt spray or other potentially harmful liquids. A specific type of joint that preferably is sealed is the joint between two structural members formed by welding or bonding with a structural adhesive. One example is in the joining of lightweight metal and plastic materials in the fabrication, repair and reconstruction of automotive and truck vehicle bodies and component panels and parts such as doors and hoods.\nA typical method for bonding such substrates in the automotive industry is described in U.S. Pat. Nos. 5,487,803 and 5,470,416, both incorporated herein by reference. According to this method, an inner panel is positioned within an outer panel with the edge of the outer panel extending beyond the edge of the inner panel. The edge of the outer panel is then folded or crimped over the edge of the inner panel in a process known as hemming and the resulting structure is known as a hem flange. In order to permanently secure the inner panel with respect to the outer panel of the hem flange a structural adhesive is applied between the surfaces of the panel or the two panels are welded together to form a permanent hem flange assembly. A sealer is applied to cover and protect the adhesive or weld joint after the joint has been formed.\nA particular problem with sealers used in connection with adhesive joints concerns cure temperature. (Meth)acrylic-based structural adhesives are becoming increasingly popular to form hem flange joints. (Meth)acrylic-based structural adhesives typically are cured at temperatures less than 121.degree. C. (250.degree. F.). Various plastisols that only cure at temperatures greater than 149.degree. C. (300.degree. F.) are currently used as hem flange sealers. However, if the previously applied (meth)acrylic-based adhesive has not itself completely cured before it is exposed during curing of the sealer to temperatures greater than 121.degree. C. (250.degree. F.), the (meth)acrylic-based adhesive tends to volatilize, thus losing its integrity. In addition, induction curing is also becoming increasing popular. Metal substrates can warp when subjected to induction curing at the temperatures required for the cure of plastisol sealants.\nSealers also must be paintable, chemically compatible with the adhesive, exhibit good flow characteristics, adhere to the substrate and joint surface and exhibit flexibility to absorb stresses caused by joint movement. Adhesion to metallic substrates can be especially problematic due to the presence of stamping, milling and drawing oils on the surface that are the result of the joint manufacturing process.\nA low temperature curable sealer that exhibits good adhesion to metallic substrates thus would be very desirable."}
{"text": "This invention relates to information storage systems and in particular to an information storage system having a receptacle wherein a recording cartridge containing a sheet of flexible information storage media is inserted for reading and writing information on the media.\nIn recording or reading information on a moving information storage media, such as a magnetic recording disc, relative motion between a magnetic head transducer used for reading and writing information and the storage media on which information is written on and read from is required. The relative velocity between the head and media interface may be, for example, 550 IPS for video and high density digital applications. To achieve maximum performance, an interface of intimate contact is provided between the head and the media surface without destruction or excessive wear of the magnetic flux responsive coating on the surface.\nAs contact pressure between the head and media is increased to improve performance, both media and head wear are increased. The problem is particularly severe in video applications and compounded with the medias utilized for storing single video frames on separate tracks of a magnetic disc wherein one track at a time is continuously in contact with the head to provide a continuous stationary display of a single frame. In only five minutes of playtime revolving at 3,600 RPM, for example, a track on the disc is scraped about 18,000 times by the head; the by-products are so hard and abrasive that the same materials are commonly used as lapping compounds.\nTo prevent failure caused by contact between the head and disc, lubricated surfaces and/or air film separations have been used. However, any separation between the head and disc caused by such lubricating fluid or air film imposes a loss of signal and hence performance. A head/disc separation equal to one wave length could cause about 54.6 db loss in the output of the replay head. Since it is desired to record wave lengths that approach 1.75 microns, the playback head voltage is reduced to 50% by only 0.19 microns of separation. On the other hand, as discussed above, reduction of separation to meet desired performance would cause the interface to be destroyed within a few seconds.\nPrior art solutions to the head-to-disc interface problem have generally been of two types: flying heads in conjunction with rigid hard-surface discs and heads having large surface areas buried in soft flexible \"floppy\" discs. Flying head discs are very expensive and require complicated and expensive recording/playback systems. A flexible or so-called \"floppy\" disc reduces the handling and cost problem incurred in the flying head, rigid-disc system. Some recording/read heads for \"floppy\" discs are relatively large to provide an interface comprised of a large contoured head buried in the soft flexible media. The large record/read head surface area distributes the force per unit area to reduce media wear and separation loss. As the media is moved past the head, however, air collects between the head and disc surface to form an air film. The thickness of this air film is a function of head and media finishes, media stiffness, head-media penetration, head size, head surface contour, viscosity of the air and disc-head relative velocity. Because of these restraints, most flexible or \"floppy\" disc applications are limited to slow speed, low bandwidth digital computer applications or voice recording systems.\nIn U.S. Pat. No. 4,175,274, means for biasing a magnetic recording disc toward the recording head by aerodynamic pressure is described and claimed. The specific embodiment described therein is a flexible spring member affixed to the inside cover of a recording cartridge. The present invention relates to an improved embodiment for biasing the disc toward the head and maintaining an intimate head-to-disc interface.\nIn recording systems known in the art the head is moved vertically as well as horizontally with respect to the media. To prevent damage to the media when the head is repositioned, the head is first moved vertically away from the media, then moved horizontally to the desired track location and finally moved vertically toward the media to assume its proper position for system operation. This procedure necessitates the use of a solenoid, servosystem or other means for moving the head both horizontally and vertically, thus adding to the number of components and the expense required to operate the system.\nMany information storage systems are adaptable for receiving portable cartridges containing the storage media therein. Such cartridges may include a tab, seal or other means which, when removed by the user, opens an electrical circuit in the information storage system to prevent information stored on the media from being erased. This allows the user to protect the information, but has the disadvantage of precluding subsequent erasure of the information on the media and use of the cartridge to record new information unless the tab or seal is somehow replaced."}
{"text": "The invention relates to encryption apparatus for audio communication through a television system. Audio security is as important as video security. With the advent of satellite links, private two-way audio channels that can be associated with a video channel are desirable, for instance for business teleconferences. Secure audio scramblers are already known for telephone lines. With video channels, however, advantage can be taken of the television system facilities to reduce the cost of scrambling the audio signal in particular by using the randomizing and synchronizing capability of video scramblers already present. Also, speech quality in the scrambling-descrambling process may be improved so as to exceed the telephone-line bandwidth. It is also important to allow for more than one audio system with the same installation.\nIt is known to scramble a digit stream as part of a video digital transmission system. See for instance \"Digital Transmission Techniques\" by G. M. Drury pp. 37-49 in IBA Technical Review (Great Britain) 1976 Vol. 9.\nIt is known from U.S. Pat. No. 3,958,077 of J. Ross et al. to generate pseudo-random numbers with digital shift registers. The patent also shows how to perform pseudo-random scanning in a television system.\nIn a scrambled television system, U.S. Pat. No. 3,717,206 of V. R. Zopf et al. shows that for subscription purposes coded signals have been associated with the transmitted video signals. This is also shown for unscrambling in U.S. Pat. No. 4,070,693 of H. B. Shutterly.\nA pseudo-random generator is disclosed in U.S. Pat. No. 3,681,708 of M. E. Olmstead.\nIt is known from U.S. Pat. No. 3,105,114 of W. Koenig to divide into segments a signal and to introduce selected time delays between such segments for scrambling effect.\nU.S. Pat. No. 3,659,046 of Angleri et al. discloses a message, in binary form, being scrambled in a pseudo-random fashion by logically combining bits.\nFrom U.S. Pat. No. 3,731,197 of J. E. Clark it is known to sample an information-bearing signal and to read the samples in a prearranged abnormal order to obtain unintelligible secured signals.\nU.S. Pat. No. 3,824,467 of R. C. French shows an audio signal divided into segments which are rearranged into a new sequence before transmission. Encoding-decoding is provided by binary pseudo-random addressing of storage devices. Each storage device transmits its stored time element while storing a new time element.\nU.S. Pat. No. 3,921,151 of G. Guanella teaches dividing an audio signal into segments of equal time intervals which are temporarily stored. Scrambling of the segments by reading-out in a random pattern is achieved so that at each location there is no repetition, nor omission.\nU.S. Pat. No. 3,773,977 of G. Guanella, which patent is related to aforementioned U.S. Pat. No. 3,921,151, shows the use of coinciding aperiodic cipher signals from which control signals are derived for determining the storage elements, together with automatic monitoring of the occupancy to avoid omissions, or repetitions.\nIt is known from U.S. Pat. No. 3,819,852 of Peter Wolf to transmit an audio signal in time-compressed form during the period of a line in the vertical blanking interval of a television system.\nThe object of the present invention is to add an audio scrambler to a television system, while using to a maximum extent the video installation for scrambling and unscrambling in the transmission and reception of an audio message."}
{"text": "The present invention is concerned with a process, and the devices for implementing it, which assures the suppression of harmful factors and permits the elimination of at least the majority of the ecologically objectionable pollutant elements--noxious or undesirable due to their toxicity, their odor, and their opaqueness--contained in the gas or liquid wastes discarded by installations manufacturing mineral fibers, and which also assures reduction of the noise produced by these same installations.\nThe invention is concerned with installations for the manufacturing of fiber blanket, mat padding, or boards of mineral fibers and especially glass, agglomerated by thermosetting or thermoplastic binders, which coat the fibers and/or bring about close binding between fibers in the finished product.\nThe binders commonly used in this type of manufacturing have a base consisting of pure or modified phenoplast or aminoplast resins, since these present advantageous characteristics for the manufacturing of agglomerated fibrous products. They are thermohardenable, soluble or emulsifiable in water, they adhere well to the fibers, and are relatively low in cost.\nGenerally, these binding agents are used dissolved or dispersed in water to which certain ingredients are added, in order to form the binder which is sprayed on the fibers.\nUnder the effect of the heat to which they are subjected during the fiber products manufacturing processes, these binders release toxic volatile elements having a perceptible pungent odor even at very weak concentrations, such as phenol, formaldehyde, urea, ammonia, and decomposition products of organic materials.\nOther binders are used for certain applications due to their very low cost. Certain extracts of natural products are hardened by drying and cross linking, such as occurs with linseed oil upon oxidation. Others are thermoplastic, as for example bitumen. During the fiber binding process, they are all, at least to some extent, increased in temperature and to a temperature sufficient to cause the release of volatile elements, noxious or otherwise undesirable, among other reasons, due to their odor.\nIn the text below, the word \"binder\" will be used to designate any one or all of the binding products mentioned above, whether they are used in liquid form, dissolved or suspended in water or in other liquid, or in an emulsion.\nThe invention relates to that part of the installation for manufacturing agglomerated fibrous materials called the fiber collecting or forming section, which is situated immediately after the fiber production apparatus, and in which the following operations are carried out primarily:\nthe conveying of the fibers from the fiber production apparatus to the mat or blanket forming equipment;\nthe application of the binder to the fibers, the binder generally incorporating pollutant elements;\nthe formation of the blanket on the fiber collecting device, for which purpose the collecting device generally consists of a perforated belt;\nthe cooling of the fibers and of the gases used for attenuating or guiding the fibers, such cooling generally being accomplished by air induced by the gases;\nthe separation of the fibers from the gases and induced air by suction these fluids through the blanket being formed; and\nthe evacuation outside of the installation of all the elements not retained by the fiber blanket or the mat being made.\nIt is in the fiber collecting or forming section that large quantities of gases and water have contact with the binder which contains the pollutant elements, and are contaminated according to a pollution process which is common to all known processes for the manufacture of blankets, mats, or boards of fibers agglomerated by a binder, and which will now be described.\n(a) The pollution of the gaseous effluents takes place according to the following process:\nThe binder is projected into the current made up of fibers and gases, coming from the fiber production apparatus, the binder being present in the form of clouds of fine droplets. Some of this binder is entrapped by the fibers, some is unavoidably deposited on the walls of the installation, and finally some is found in the gases or fumes in the form of fine droplets and in the form of vapor.\nThus two fluid contamination modes coexist, the one consisting of contamination by droplets of the binder and the other consisting of contamination by vapors of the binder. In the binder application, the binder atomization or dispersion devices used furnish particles or droplets within a very wide range of diameters. The finest droplets are not entrapped by the fibers and are drawn through the blanket being formed by the gaseous current, in which they are present in suspension.\nThe droplets of binder deposited on the fibers during the binder application are subjected to the kinetic effects of the gaseous current passing through the blanket being formed. A large quantity of droplets is extracted from the fibers, migrates through the blanket, and is found in suspension in the exhausted gases.\nFinally, the desire to obtain a homogeneous distribution of the binder in the blanket makes it necessary to disperse the binder in the fiber and gaseous current in an area situated near the fiber production apparatus, where this current still has a well-defined geometric form but where its temperature may still be high enough so that some of the binder, or at least its most volatile components, are evaporated. These pollutant vapors mix with the gases and contaminate them.\nIn the text below, the word \"fumes\" will be used to designate the gaseous effluents which pass through the fiber blanket and are evacuated outside of the collecting unit, i.e., the gases used for attenuating or guiding the fibers, the fluids induced by these gases, and the pollutant elements in the form of droplets or vapor suspended in these fluids. It is to be understood that various features of the invention, such as treatment steps and components of the apparatus, may be employed with \"fumes\" having a wide range of compositions and pollutants. It is preferred to treat all components of such fumes, but various features of the invention may also be employed with gases originating in fiber production operations in which the gases have pollutant components, whether or not the pollutants have their origins in fiber binders.\n(b) The functions performed by the water in a fiber collecting unit make a large degree of pollution inevitable in any installation in which binders are used.\nIn operation, water is used:\n(1) to dilute and carry the binder when the latter is used in liquid form;\n(2) to wash or scrub the fumes, an operation which consists:\n(2a) of causing the largest possible amount of pollutants contained in the fumes in the form of droplets or vapor to be captured by the droplets of the scrubbing water, thus causing the pollutant charge of the fumes to be transferred to the wash water; PA1 (2b) of capturing and entraining on the walls of the collecting unit the fibers suspended in the fumes;\n(3) to wash the different parts of the collecting installation (perforated belt, fume flues, etc.) in order to evacuate the binder and the fibers deposited therein.\nDuring these operations the wash water is charged with binder components which are soluble, insoluble, or in the vapor state, and the concentration of pollutant elements may reach high values.\nThe foregoing description of the manner in which the fumes and the water are contaminated is based on an interpretation of measurements and observations made in actual manufacturing installations.\nData derived from such measurement and observation is herein given by way of information; but it will be understood that other data and explanations may be found and that the invention is not limited by the data given.\nIn all installations for the manufacturing of agglomerated fibrous products, regardless of the fiberization process used, the effluent pollution described above involves considerable quantities of effluents.\nIn installations equipped with devices for attenuating fibers by blowing, in which the material to be attenuated is transformed into fibers by means of high-energy jets, the quantities of fumes discharged into the atmosphere are--for the best known processes--on the order of magnitude of the following values:\n100 Nm.sup.3 per kilo of fibers for the process described in the Slayter U.S. Pat. No. 2,133,236;\n300 Nm.sup.3 per kilo of fibers for the AEROCOR process (Stalego U.S. Pat. No. 2,489,243);\n70 Nm.sup.3 per kilo of fibers for the SUPERTEL process (Levecque French Pat. No. 1,124,489 and U.S. Pat. Nos. 3,114,618 and 3,285,723);\nwhich, for large production plants, leads to outputs ranging from 500,000 to 1,000,000 Nm.sup.3 /hr. (In these values, Nm.sup.3 refers to the cubic meter volume at standard atmospheric pressure and room temperature.)\nIn installations equipped with fiber attenuating devices, in which the material to be attenuated is transformed into fibers under the effect of mechanical forces--centrifugal for example--and where a gaseous current is only used as a medium (generally flowing in an essentially horizontal direction, see FIG. 14, for example) for carrying the fibers produced towards the collecting device--the quantity of fumes given off is a little less, but nevertheless very important: for example, 30 Nm3 per kilo of fibers, for the process described in Powell U.S. Pat. No. 2,577,431, which for a production plant results in outputs on the order of 300,000 to 400,000 Nm.sup.3 /hr.\nThe quantities of polluted water are pretty much the same for all processes, and on the order of 1,000 m.sup.3 /hr. or more for large industrial installations.\nThe volume of these quantities of polluted effluent has led legislatures first to limit the concentration of phenol compounds in the effluents discarded in the atmosphere, and later to prohibit discarding of any pollutants, at least in certain countries.\nFurthermore, limitations concerning the odors or the opaqueness of discharged effluents have been established in various countries.\nIn addition, installations for the manufacture of agglomerated fibrous products also tend to pollute in another respect. In addition to toxic or pungent-smelling products, these installations discard substantial quantities of steam, on the order of 20 to 30 metric tons per hour for large plants, which steam escapes from stacks in very opaque plumes.\nNoise is another type of nuisance created by installations for the manufacturing of agglomerated fibrous products. In these installations, the noise is essentially emitted by two sound sources--the apparatus for producing the fibers and the fan for extracting the fumes.\nActually, all the equipment for producing fibers mounted in these installations uses jets of gases at high speed either for transforming the material to be drawn or attenuated into fibers or for directing the fibers produced. It is known that the acoustic power level emitted by these jets considerably increases with the speed of the jets. This level may exceed 100 decibels adjacent to the fiber production apparatus, where the operators are required to work. This level is much higher than the level tolerated by industrial regulations in many countries.\nFurthermore, the acoustic power developed by the fume extraction fan is transmitted along the flues connecting with the fume exhaust stack. The latter is ordinarily situated outside the buildings, where it functions as an antenna, and radiates this acoustic power into the surrounding environment. The inconvenience resulting for the vicinity has caused authorities in different countries to order the shut-down of certain installations.\nThe need to reduce or eliminate the pollution produced--and this at costs low enough not to overly influence the cost price of the finished product--is pressing. Numerous investigations have been carried out on this problem, and certain solutions have been developed.\nThe process according to the invention is characterized by the fact that the fumes (as hereinabove defined) are partially recycled, so as to cause them repeatedly to traverse the blanket or mat being formed. The process according to the invention is also characterized by the fact that the majority of the heat contributed by the gases coming from the fiber production apparatus and the attenuated fibrous material is transferred to the wash water, by the fact that the wash water is cooled, by the fact that the fumes are washed in water after they have traversed the blanket or mat and the fiber collecting device in order to transfer to the water some of the pollutant products contained in these fumes, by the fact that the non-recycled part of the said fumes is purified before evacuation into the atmosphere, by the fact that at least some of the wash water is recycled--a certain quantity of which has been subjected to a treatment for extracting at least a sizable fraction of the pollutant products contained in the wash water--and by the fact that the solid wastes are subjected to a purification treatment before final disposal.\nThe foregoing process effects cooling of the recirculation gases, which is important in making possible such recirculation. In combination with such cooling of the recirculating gases, it is preferred also to spray water on the current of fibers and gases in the receiving chamber, in order to cool the fiber and gas current. Such water spraying, with resultant cooling of the current, together with the cooling of the recirculating gases provides for reduction of the temperatures of the current notwithstanding the substantial absence of induction of ambient air by the attenuating blast.\nAccording to a particularly important characteristic of the invention, the quantity of fumes discarded into the atmosphere is essentially equal to the quantity of gases flowing from the attenuating device.\nThe invention is particularly concerned with recycling the majority of the fumes in the installation, and with treating and evacuating only a small portion of the fumes--it being possible for the recycled portion to reach at least 95% of the total quantity of fumes ordinarily evacuated into the atmosphere. The quantity of fumes to be purified before discarding may thus be less than 5% of all of the fumes, which even makes it practicable to use costly purification treatment, whose effectiveness is total--as for example burning--without prohibitive energy expenditures.\nAnother object of the invention is to render insoluble the thermohardenable resins contained in the water. These resins are rendered insoluble, according to the invention, by means of a heat treatment--preferably at a temperature greater than 100.degree. C., and more advantageously ranging between approximately 150.degree. and 240.degree. C., and under pressure.\nThe application of the above process (for rendering resins insoluble) to at least some of the cooling and washing water is advantageously used to render insoluble the dissolved binder components contained in the water, in order to subsequently be able--by means of known techniques--to extract insoluble materials and thus to maintain the concentrations of the pollutant constituents in the washing and cooling waters at a level compatible with the continuous re-utilization of these waters in the installation. The wash water thus circulates in a closed circuit and any external rejection of pollutants with the wash water is eliminated.\nAnother object of the invention consists of a heat treatment to which the wash water is subjected--a treatment which consists of vaporizing it and of heating this vapor to a temperature sufficient so that the pollutant constituents are transformed into non-pollutant constituents.\nThe invention also is concerned with means for sound insulation--adjusted to the particular configurations shown--to the devices for conveying and guiding the recycled fumes, in order to reduce the noise emitted by these devices, and with a particular arrangement of the apparatus for evacuating the non-recycled fumes into the atmosphere, which reduces the noise emitted by this apparatus in the surrounding environment.\nIn addition to the general objectives above referred to the present invention also contemplates certain controls for the operating conditions, as pointed out just below.\nIn techniques of the kind briefly referred to above and described in detail hereinafter, the use of the various means for suppression of pollution, especially the recirculation of the current of the attenuating gas and also the separation of the pollutants from the recirculating gas, as by means of a water spray, may at times tend to introduce undesirable fluctuations in the conditions under which the fibers are formed or attenuated, and the condition under which the fiber blanket is formed. Because of the recirculation of a large part of the gases, it is desirable to more completely enclose the forming section, than has been customary where the suppression of pollution by recirculation of the gases is not contemplated. With the more tightly enclosed forming section and where recirculation of gases is employed for the purpose of suppression of pollution, there may be tendencies for fluctuation of both the pressure and the temperature of the gas in the forming section. The pressure will vary in accordance with the quantity of the gases which are diverted and discharged from the recirculation flow path; and in addition, the temperature will vary in accordance with a number of factors including not only the quantity of gas diversion and discharge from the recirculation flow path, but also the extent of water spraying utilized for separation of pollutants from the recirculation gases, as well as the temperature of the water used for such water spraying. Still further, variation in atmospheric conditions, for example as between summer and winter, may also influence the operating conditions with respect to both pressure and temperature.\nVariable factors such as those just referred to tend to alter uniformity of fiber and fiber blanket production, particularly in the fiber formation by gas blast attenuation, since uniformity of the fibers depends in part upon uniformity of the conditions of temperature and pressure. In fact, if the temperature of the gaseous current and consequently of the fiber blanket is too high, polymerization of the binder will start prematurely, i.e., in for forming section, instead of awaiting feed of the blanket into the binder curing oven. This condition tends to reduce the mechanical properties of the products, particular their resilience.\nOn the other hand if the temperature of the gases and consequently that of the blanket is too low, the moisture carried by the blanket increases, and this reduces the efficiency of the curing oven, and can even lead to dimensional irregularities of the manufactured products.\nPressure variations tend to adversely influence the effectiveness of the devices used to reduce the pollution in the gases discharged through the stack. A negative pressure in the formation chamber, that is a pressure below atmospheric pressure will increase the quantity of the air penetrating into the forming section and consequently the quantity of gases to be diverted from the recirculation path and discharged. This results in an increase in the quantity of pollutants ejected into the atmosphere. A positive pressure, on the other hand, leads to leakage or discharge from the formation chamber of gases not yet treated, thereby impairing the intended suppression of pollution.\nWith the foregoing in mind it is contemplated according to the present invention that controls be provided for maintaining substantial uniformity of the conditions prevailing in the zones of fiber attenuation and fiber blanket formation, particularly uniformity of pressure and temperature of the gases in these zones. In addition, it is further contemplated to regulate the volume of the gas in circulation.\nIt is also contemplated according to the present invention that the controls for temperature and pressure be adjustable in order to establish the desired pressure and temperature levels.\nOther objects and advantages of the invention, including in particular numerous specific advantages for the recycling of fumes, will be given and explained more completely below."}
{"text": "Social media platforms are becoming an increasingly important way to interact and engage with people, such as customers of a business. Both customers and businesses can benefit from an increase in the number of opportunities to engage through the use of social media. In particular, the benefits of engaging are particularly pronounced when a customer attends an event at the physical, geographic location of a business or event. The ability to recognize the customer's physical presence and engage with the customer in real time, can be used to enhance the customer's experience at the location, and thus is extremely valuable.\nHowever, some techniques for determining the presence of a customer at a physical location using social media are extremely limited in their ability to effectively determine relevant location information, particularly in real time. For example, some existing techniques rely solely on social media posts with geotagged physical locations in order to determine social media activity by customers that are physically present at a location. However, techniques that rely solely on geotagged social media posts can miss the vast majority (e.g., up to about 95%) of social media activity relating to the geographic location of interest."}
{"text": "Digital graphic design, image editing, audio editing, and video editing applications (hereafter collectively referred to as media content editing applications or media editing applications) provide graphical designers, media artists, and other users with the necessary tools to create a variety of media content. Examples of such applications include Final Cut Pro® and iMovie®, both sold by Apple, Inc. These applications give users the ability to edit, combine, transition, overlay, and piece together different media content in a variety of manners to create a resulting media project. The resulting media project specifies a particular sequenced composition of any number of text, audio, image, and/or video content elements that is used to create a media presentation.\nVarious media editing applications facilitate such composition through electronic means. Specifically, a computer or other electronic device with a processor and computer readable storage medium executes the media content editing application. In so doing, the computer generates a graphical interface whereby designers digitally manipulate graphical representations of the media content to produce a desired result.\nOne difficulty in media editing is that a user cannot conveniently access various candidate clips for a particular section of a presentation. Typically, if a user wishes to evaluate different clips than those included in the composite presentation, the user must select from a large group of clips by iteratively accessing various media libraries.\nAnother difficulty in media editing is that a user cannot conveniently manage candidate clips (i.e., clips that are not included in a composite presentation) such that the clips can be easily accessed when modifying a particular section of the composite presentation. At any given time, the user may wish to associate one or more candidate clips (e.g., alternative “takes” of a scene, alternative content, etc.) with a particular section of a presentation. Thus, there is a need for a media editing application that allows a user to associate candidate clips with various sections of a composite presentation, thus allowing the user to retrieve the associated clips when evaluating a particular section of the composite presentation."}
{"text": "Tires, particularly commercial vehicle tires as used on trucks, must be periodically removed from service and replaced as the tires eventually wear from use. Depending upon the application and maintenance of each tire, anomalies in the tread region can also develop during use that lead to the tire being removed from service before reaching the normal useful life of the tread. Typically, such anomalies are discovered as the driver senses a change in the ride of the vehicle or a change in the tire noise during vehicle use, and removal of a tire having anomalies is frequently a subjective determination. Such anomalies can include, for example, depressions or uneven wear in the tread region, which is generally caused by unequal stress distribution laterally across the tread region. “Heel and toe” wear is a form of irregular wear that can occur in which one side of a tread element—such as the leading edge of a tread block—wears more quickly than the opposing side. An explanation of various types of irregular tire wear can be found in a guide entitled “Radial Tire Conditions Analysis Guide: A Comprehensive Review of Tread Wear and Conditions”, published by The Truck Maintenance Council (TMC).\nMany different factors can contribute to irregular tire wear. For example, the operating conditions of the vehicle as well as the type of load carried can both have an impact on tire wear. In addition to the existence of many possible causes, it is common that a combination rather than a single factor will start uneven wear or cause it to increase. Determining a solution where a combination of factors is contributing can be particularly difficult.\nAccordingly, a need remains to reduce the irregular wear occurring on tread elements such as e.g., tread blocks. The present invention provides one or more features than can help resist the onset and continuation of irregular wear. The onset of “heel to toe” wear is a particular concern that the present invention can be useful in addressing."}
{"text": "A long term evolution advanced (LTE-Advanced) network standard has been developed to provide wireless data rates of 1 Gbps downlink and 500 Mbps uplink. The LTE-Advanced network standard also offers multi-carrier transmission and reception within a single band as well as multi-carrier transmission and reception within two separate bands. Multi-carrier transmission within a single band is referred to as intra-band transmission and reception. In contrast, multi-carrier transmission and reception within two different bands is referred to as inter-band transmission and reception. LTE-Advanced technology is also known as fourth generation (4G) technology.\nMajor challenges are associated with realizing front end radio architectures (FERAs) that are necessary for providing multi-carrier operation using intra-band and inter-band transmission and reception. One of the major challenges pertains to transmit (TX) multi-carrier inter-modulation (IMD) products that can fall within a multi-carrier receive (RX) channel. Another major challenge pertains to RX de-sense due to the proximity of a multi-carrier TX channel. This major challenge is exacerbated by multi-carrier operation in which an offset frequency between the TX channel and the RX channel can be lower than a typical duplex frequency offset. Therefore, a need exists for a FERA having improved IMD performance."}
{"text": "1. Technical Field\nThe present disclosure relates to energy plant optimization and, more specifically, to a system and method for energy plant optimization using mixed integer-linear programming.\n2. Discussion of Related Art\nLarge-scale commercial buildings such as skyscrapers and industrial buildings such as factories have large and complex energy generation and utilization needs. These needs include chilled water, chilled air, hot water, hot air, electrical loads, and the like. To meet these complex energy generation and utilization needs, these buildings, referred to herein as energy plants or simply plants, may include a large number of devices that utilize energy, generate energy, and/or store energy, for example, to provide for the plant's heating, cooling and electricity needs. These devices are referred to herein as energy devices and include devices that utilize energy such as water chillers, air conditioners, water heaters, and devices that utilize plant electricity. Energy devices may also include devices that generate energy such as fossil fuel-powered electrical generators, geothermal systems, solar collectors and photovoltaic cells. Energy devices may also include devices that store energy for later use such as thermal energy storage devices, batteries and the like.\nConventionally, energy systems may be activated and deactivated as needed. For example, a water chiller may activate when there is a need for chilled water and an air conditioner may activate when there is a need for cool air. However, as a plant may include a great number of energy devices, it is possible and in fact probable that many such devices expend energy to work against each other and operate in a less than ideal fashion.\nSuboptimal plant energy generation, storage and usage may lead to excessive costs, unnecessarily large environmental impact, overtaxing of energy grids, and difficulty in complying with present and proposed governmental regulations."}
{"text": "This is a continuation of application Ser. No. 819,099, filed Mar. 4, 1997, now U.S. Pat. No. Re. 37,058, issued Feb. 20, 2001, which is a continuation of application Ser. No. 360,226, filed Dec. 20, 1994, now abandoned, which is a broadening reissue application of U.S. Pat. No. 5,173,814, issued Dec. 22, 1992 from application Ser. No. 653,100, filed Feb. 8, 1991, said application Ser. No. 653,100 being is a continuation of application Ser. No. 07/402,917, filed Sep. 5, 1989, now U.S. Pat. No. 5,001,581, issued Mar. 19, 1991, which is a continuation of application Ser. No. 201,736, filed Jun. 2, 1988, now  U.S. Pat. No. 4,894,738, issued Jan. 16, 1990, now U.S. Pat. No. Re. 35,792, issued May 12, 1998, which is a continuation-in-part of application Ser. No. 038,049, filed Apr. 14, 1987, now U.S. Pat. No. 4,843,500, issued Jun. 27, 1989, which is a continuation-in-part of application Ser. No. 767,671, filed Aug. 21, 1985, now U.S. Pat. No. 4,658,312, issued Apr. 14, 1987, which is a continuation of application Ser. No. 412,093, filed Aug. 27, 1982, now abandoned, which is a continuation-in-part of application Ser. No. 326,559, filed Dec. 2, 1981, now U.S. Pat. No. 4,519,010, issued May 21, 1985, said application Ser. No. 412,093 also being a continuation-in-part of application Ser. No. 244,971, filed Mar. 18, 1981, now abandoned, said application Ser. No. 201,736 also being a continuation-in-part of application Ser. No. 32,954, filed Mar. 31, 1987, U.S. Pat. No. 4,779,165, issued Oct. 18, 1988, now U.S. Pat. No. Re. 34,412, issued Oct. 19, 1993, which is a continuation of application Ser. No. 733,231, filed May 10, 1985, now abandoned, which is a continuation-in-part of the said application Ser. No. 412,093."}
{"text": "1. Field of the Invention\nThe present invention relates to computer networks and in particular to handling client requests in systems having multiple servers.\n2. Description of Related Art\nOn computer networks, nodes may interact through requests for data, services, or other resources. Client nodes may generate requests and server nodes may service those requests. Nodes may be stand-alone computers, servers, or other computing devices, as well as virtual machines, threads, processes, or other network participants. To handle a greater volume of requests, a plurality of servers may be allocated to responding to one or more types of requests and the servers as a group may be referred to as a server farm. For example, a bank, ticket office, or some other organization may use a server farm to process requests from clients. The servers may be web servers on the Internet responding to web clients.\nTo provide coherent service to clients, state or session information may be collected and used by servers, so that their response to a given client may depend on session data for the client in the server farm. For example, session information may be used to implement a “shopping cart” used across multiple requests over one or more servers. It may be more efficient to store the session information on a single server within the server farm, rather than disperse that information, and ensure that the corresponding server services the requests from the associated client.\nA typical system may have a unique session ID assigned to correspond to collected client session information. Initially, requests do not contain session IDs since there is not yet a session formed with the server farm. The first server to receive a request from a given client typically generates the session I). After a server creates a session for a client, the server generates a session ID to uniquely represent that session, and instructs the client to return the generated session ID in subsequent requests from that client. For example, a session ID may be encoded in the URLs the client is instructed to access or it may be stored in a server-generated cookie the client is instructed to send back to the server with each request.\nRequests may be routed between the clients and the servers through one or more devices, which may be considered upstream from the servers and downstream from the clients. These devices may be part of the server farm or may operate within the network or other system. For example, the server farm may include one or more routers that receive the requests from the network and route them to servers.\nThe work that a server performs processing requests it receives is referred to as the load of that server. Depending on the allocation of requests and the processing required, the load may differ widely between servers and limit performance. Another example of upstream devices is load balancers. In order to distribute load across a plurality of servers within a server farm, load balancers may be used. Load balancers may be additional devices, may be part of other devices, or may include other functionality. For example, load balancers may receive incoming requests to the server farm and route the requests to individual servers in addition to determining where requests may be sent to better balance loads in the server farm.\nIn some systems, a load balancer or other device may be configured to receive incoming requests, examine them for a session ID, and then route them to the server that generated the session ID so that session information is kept on a single server within the server farm. The technique of using the session ID to select the server that will service a given request is called “sticky sessions”, “sticky load balancing”, or “sticky routing” since the presence of a session ID causes the requests associated with that session to “stick” to a specific server.\nSome systems may include load balancers or other devices in the server farm that receive requests and send them to servers. These devices may maintain a mapping of session IDs to servers. These devices may then compare the session IDs with the requests to their mapping information to find the corresponding server. However, the size of the data for the mapping may be prohibitive as there may be many session IDs-at least one per client of the server farm.\nIn some systems, servers may provide their address, such as an Internet hostname or address, to clients so that the client or devices within the server farm, such as load balancers, may use that information to route subsequent requests using standard Internet routing protocols to the corresponding server. However, exposing the server addresses or other details of the sever farm configuration communicates information about the internal arrangement of the server farm to the client or others outside the server farm. As the client is typically a distrusted third party, exposing information about the internal operation of the server farm may be a security risk.\nAlternatively, an administrator may configure a load balancer and servers with a particular ID corresponding to each server as identifying information. Servers may then include this identifying information in their session ID without revealing information about the internal structure of the server farm. However, the load balancer and server must then share data that allows the load balancer to target a specific server from the identifying information. This implementation imposes a burden on server farm administrators, as configuration information for these session IDs is then maintained for all load balancers and servers."}
{"text": "a. Field of the Invention\nThe field of the invention pertains to manifolds for splitting a fluid flow, more specifically, manifolds that uniformly split fluid from an inlet passage among a multiplicity of outlet passages.\nb. Background Art\nSplitting a quantity of fluid into equal portions is useful for many reasons. For example, a sample of fluid may be evenly distributed among the wells of a microtitration plate (\"plate\" or \"microplate\") or among a group of test tubes. In this manner, research may be conducted on several uniform samples at the same time, accelerating the research and allowing for direct correlation of results.\nOne existing method used to split a quantity of fluid into approximately equal portions involves dedicating a costly pump or syringe to each outlet passage. A common and less expensive method used for dispensing fluid into a microplate entails passing the fluid through a manifold. The instant invention relates to this latter method.\nCustomarily, a manifold comprises multiple ports or passages inserted into an elongated cylindrical cavity. Fluid is introduced into the cavity and subsequently dispersed through outlet ports into a receptacle (e.g., a microplate). A manifold with only two symmetrically placed outlet passages would have a nearly equal flow from its outlet passages. However, manifolds with more than two outlet passages are the most useful, and the more outlet passages, the greater the flow disparity.\nIn a typical manifold, the number of outlet passages is designed to match the number of wells in a single column of a microplate, which has several columns of wells. In this manner, the manifold may fill the microplate one column at a time. Thus, Manifolds having either eight or twelve outlet passages are common. The outlet passages are positioned, not necessarily symmetrically, on each side of the inlet passage. Accordingly, the flow from each of the outlet passages is not uniform in the absence of a flow modifying device such as the instant invention.\nIn the operation of a typical manifold for microplate dispersion applications, the manifold is first primed to eliminate the manifold's \"dead volume\" (i.e., the volume of fluid needed to fill the manifold before useful dispensing can start). Then, a volume of fluid equal to the volume desired for each well in a single column of the microplate times the number of outlet passages (which equals the number of wells in a single column of wells), is dispensed into the inlet passage of the manifold. Ideally, the manifold then evenly disperses the flow among the microplate wells of the column. However, the flow from the outlet passages nearest the inlet passage is typically less than the flow from the outlet passages most distant from the inlet passage. This discrepant flow pattern results from the internal flow dynamics of the manifold, which must be adjusted to balance the flow."}
{"text": "The technology is known that controls a welding electric current in a high precision by performing a primary electric current control through an inverter in a welding transformer of resistance welding device (referring to a Patent document 1). Further, through devising its coil structure, a further high speed control is attempted (referring to a Patent document 1).\nResistance welding technology is widely used for the manufacture line of a car or general industries. In those industries, in order to overcome competition on a world level, improvement in productivity is indispensable.\nFurthermore, development of the energy-saving technology which cuts down CO2 from the viewpoint of earth-environment protection is also an urgent subject. However, the conventional resistance welding method for the steel plate with bad weld ability was presenting the aspect which moved against productivity or energy saving. The various resistance welding methods that it was suitable for such various uses are developed (referring to a Patent document 1,2,3,4)."}
{"text": "In the papermaking industry, it is generally known that paper to be converted into a consumer product such as paper towels, bath tissue, facial tissue, and the like is initially manufactured and wound into large rolls. By way of example only, these rolls, commonly known as parent rolls, may be on the order of 10 feet in diameter and 100 inches across and generally comprise a suitable paper wound on a core. In the usual case, a paper converting facility will have on hand a sufficient inventory of parent rolls to be able to meet the expected demand for the paper conversion as the paper product(s) are being manufactured.\nBecause of the soft nature of the paper used to manufacture paper towels, bath tissue, facial tissue, and the like, it is common for parent rolls to become out-of-round. Not only the soft nature of the paper, but also the physical size of the parent rolls, the length of time during which the parent rolls are stored, and the fact that roll grabbers used to transport parent rolls grab them about their circumference can contribute to this problem. As a result, by the time many parent rolls are placed on an unwind stand they have changed from the desired cylindrical shape to an out-of-round shape.\nIn extreme cases, the parent rolls can become oblong or generally egg-shaped. But, even when the parent roll is are only slightly out-of-round, there are considerable problems. In an ideal case with a perfectly round parent roll, the feed rate of a web material coming off of a rotating parent roll can be equal to the driving speed of a surface driven parent roll. However, with an out-of-round parent roll the feed rate can likely vary from the driving speed of a surface drive parent roll depending upon the radius at the web takeoff point at any moment in time.\nWith regard to the foregoing, it will be appreciated that the described condition assumes that the rotational speed of the parent roll remains substantially constant throughout any particular rotational cycle of the parent roll.\nIf the rotational speed remains substantially constant, the feed rate of a web material coming off of an out-of-round parent roll will necessarily vary during any particular rotational cycle depending upon the degree to which the parent roll is out-of-round. In practice, however, parent rolls are surface driven which means that if the radius at the drive point changes, the rotational speed can also change generally causing variations in the feed rate. Since the paper converting equipment downstream of the unwind stand is generally designed to operate based upon the assumption that the feed rate of a web material coming off of a rotating parent roll will always be equal to the driving speed of the parent roll, there are problems created by web tension spikes and slackening.\nWhile a tension control system is typically associated with the equipment used in a paper converting facility, the rotational speed and the takeoff point radius can be constantly changing in nearly every case. At least to some extent, this change is unaccounted for by typical tension control systems. It can be dependent upon the degree to which the parent roll is out-of-round and can result in web feed rate variations and corresponding tension spikes and slackening.\nWith an out-of-round parent roll, the instantaneous feed rate of the web material can be dependent upon the relationship at any point in time of the radius at the drive point and the radius at the web takeoff point. Generally and theoretically, where the out-of-round parent roll is generally oblong or egg-shaped, there will be two generally diametrically opposed points where the radius of the roll is greatest. These two points will be spaced approximately 90° from the corresponding generally diametrically opposed points where the radius of a roll is smallest. However, it is known that out-of-round parent rolls may not be perfectly oblong or elliptical but, rather, they may assume a somewhat flattened condition resembling a flat tire, or an oblong or egg-shape, or any other out-of-round shape depending upon many different factors.\nRegardless of the exact shape of the parent roll, at least one point in the rotation of the parent roll exists where the relationship between the web take off point radius and the parent roll drive point radius that results in the minimum feed rate of paper to the line. At this point, the web tension can spike since the feed rate of the web material is at a minimum and less than what is expected by the paper converting equipment downstream of the unwind stand. Similarly, there can exist at least one point in the rotation of the parent roll where the relationship between the web take off point radius and the parent roll drive point radius results in the maximum feed rate of paper to the line. At this point, the web tension can slacken since the feed rate of the web material can be at a maximum and more than what is expected by the paper converting equipment downstream of the unwind stand. Since neither condition is conducive to efficiently operating paper converting equipment for manufacturing paper products such as paper towels, bath tissue and the like, and a spike in the web tension can even result in a break in the web material requiring a paper converting line to be shut down, there clearly is a need to overcome this problem.\nIn particular, the fact that out-of-round parent rolls create variable web feed rates and corresponding web tension spikes and web tension slackening has required that the unwind stand and associated paper converting equipment operating downstream thereof be run at a slower speed in many instances thereby creating an adverse impact on manufacturing efficiency.\nWhile various efforts have been made in the past to overcome one or more of the foregoing problems with out-of-round parent rolls, there has remained a need to successfully address the problems presented by web feed rate variations and corresponding web tension spikes and web tension slackening."}
{"text": "Intraluminal, intracavity, intravascular, and intracardiac treatments and diagnosis of medical conditions utilizing minimally invasive procedures are effective tools in many areas of medical practice. These procedures are typically performed using imaging and treatment catheters that are inserted percutaneously into the body and into an accessible vessel of the vascular system at a site remote from the vessel or organ to be diagnosed and/or treated, such as the femoral artery. The catheter is then advanced through the vessels of the vascular system to the region of the body to be treated. The catheter may be equipped with an imaging device, typically an ultrasound imaging device, which is used to locate and diagnose a diseased portion of the body, such as a stenosed region of an artery. For example, U.S. Pat. No. 5,368,035, issued to Hamm et al., the disclosure of which is incorporated herein by reference, describes a catheter having an intravascular ultrasound imaging transducer.\nFIG. 1a shows an example of an imaging transducer assembly 1 known in the art. The imaging transducer 1 is typically within the lumen 10 of a guidewire (partially shown), having an outer tubular wall member 5. To obtain an image of a blood vessel, the imaging transducer assembly 1 may be inserted into the vessel. The transducer assembly 1 may then rotate while simultaneously emitting energy pulses, e.g., ultrasound waves, at portions of the vessel from within the vessel and receiving echo or reflected signals.\nTurning to FIG. 1b, it is known in the art that an imaging console 20 having a display screen, a processor and associated graphics hardware (not shown) may be coupled with the imaging transducer assembly 1 to form a medical imaging system 30. The imaging console 20 processes the received echo signals from the imaging transducer assembly 1 and forms images of the area being imaged. To form the images, the imaging console 20 draws multiple lines, known as “radial lines”, (not shown) on the display screen that each correspond to an angular position of the transducer assembly 1. The processor of the imaging console 20 assigns brightness values to pixels of the lines based on magnitude levels of the echo signals received from the transducer assembly 1 at the angular positions corresponding to the lines. A drawing that includes a large number of these radial lines results in an image such as an intravascular ultrasound (IVUS) image (not shown). Such an image may show, among other things, the texture of the area being imaged, such as the smoothness or the roughness of the surface of the area being imaged.\nAn example of an image 70 having a large range of magnitudes and a number of texturally distinct regions 80 is shown in FIG. 1c. Texture and the correct discrimination of the underlying surface are important in medical imaging. Such information is helpful to radiologists and other clinicians who seek to diagnose pathology. It is often the case in medical imagery that an abnormality is detectable only as a subtle variation in texture. Accordingly, an improved system and method for detecting and presenting such textural information would be desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to a common rail that is mounted to a pressure accumulation fuel injection system and accumulates high-pressure fuel.\n2. Description of Related Art\nA common rail 101 of a related art will be described with reference to FIG. 14. The common rail 101 is a pressure accumulator for accumulating high-pressure fuel pumped from a high-pressure fuel pump (supply pump or the like) and has an accumulation chamber (central hole) 123 formed therein for accumulating the high-pressure fuel therein. The common rail 101 has pipe joints 121 to be connected with external pipes (a high-pressure pump pipe, an injector pipe, and the like). A center portion of an outer end of the pipe joint 121 communicates with the accumulation chamber 123 through an inside-outside communication hole 124 (branch hole).\nAn orifice 130 for reducing pressure pulsation caused by an injection operation of an injector or an orifice 130 for reducing pressure pulsation caused by pumping operation of the high-pressure fuel pump is formed in the inside-outside connection hole 124. The orifice 130 of the related art is formed by directly drilling a main body 120 (hereinafter referred to as a rail main body) of the common rail 101. The orifice 130 is formed in the deepest portion of the inside-outside connection opening 124 because of constraints concerning the drilling. That is, as shown in FIG. 14, the orifice 130 opens in the accumulation chamber 123.\nHigh-pressure fuel is accumulated in the accumulation chamber 123 and therefore high pressure is applied to the inner peripheral surface of the accumulation chamber 123. The orifice 130 having a small diameter opens in and crosses the inner peripheral surface of the accumulation chamber 123. A portion of the opening that crosses the inner peripheral surface is hereinafter referred to as a crossing hole. As the crossing hole becomes smaller, a larger stress is concentrated at the edge of the opening of the crossing hole. Therefore, a common rail, in which the orifice 130 is integrally formed in the rail main body 120 by the drilling, is used for a pressure accumulation fuel injection system of a small size vehicle or the like, in which accumulation pressure of the accumulation chamber 23 is comparatively low (for example, 180 MPa or under).\nIn recent years, it has been required to increase the common rail pressure to 180 MPa or over also in the small size vehicle for the purpose of improving exhaust performance and the like. However, a common rail having the orifice 130 integrally formed in the rail main body 120 by the drilling has the small crossing hole, which is formed by the orifice 130, and therefore has difficulty in securing a safety factor concerning fatigue strength. Therefore, not only in a large size vehicle but also in the small sized vehicle, for the purpose of securing the safety factor concerning the fatigue strength, as shown in FIG. 3, there has been proposed a common rail 1, in which an orifice 30 is formed in a bush 31 separate from a rail main body 20 and the bush 31 is fixed in an inside-outside communication hole 24 to enlarge the crossing hole (for example, as described in Patent document 1: JP-A-2002-322965) in place of forming the orifice 130 directly in the rail main body 120.\nA technology of the related art that fixes the bush 31 having the orifice 30 in the inside-outside communication hole 24 press-fits an outer peripheral surface of the bush 31 into the inside-outside communication hole 24. When the bush 31 is inserted into the inside-outside communication hole 24, there is a possibility that a press-fitting burr A is produced at an end portion (hereinafter referred to as a press-fitting end portion) of the bush 31 opposite to the press-fitting direction of the bush 31 as shown in FIG. 4A.\nWhen the engine is driven in a state where the press-fitting burr A is produced as shown in FIG. 4A, there is a possibility that the press-fitting burr A is dropped by an action of the flow of the fuel and flows as a foreign matter to the injector. When the press-fitting burr A is dropped and flows to the injector, the press-fitting burr A clogs a filter disposed in the injector. When a small press-fitting burr A passes through the mesh of the filter of the injector, there is a possibility that the burr A enters the injector having high accuracy and causes malfunctions such as faulty sliding and faulty sealing."}
{"text": "In general, an ink jet recording head for use in an ink jet recording apparatus is composed of an ink jet recording head for forming droplets of a liquid such as ink, and of a supply system for supplying the ink and the like to the head.\nWith regard to a connection of a wiring board to a liquid discharge substrate, which are applied to the recording head, a wide-array ink jet apparatus composed of a printhead substrate in which an electrical connection electrode is present on a surface opposite to a surface having a discharge port therein has been disclosed in JP 11-192705A. In FIG. 18 and FIG. 19, a wide-array ink jet pen 210 disclosed in JP 11-192705A is shown. FIG. 18 is a perspective view of the wide-array ink jet pen having a wide-array printhead. FIG. 19 is a cross-sectional view of a part including a printhead die and a carrier substrate 220, showing electrical connection portions of the wide-array ink jet printhead of FIG. 18. The pen 210 is composed of a wide-array printhead 212 and a pen body 214. The pen body 214 is a housing to which the printhead 212 is attached. In the pen body 214, an internal chamber 216 serving as a local ink reservoir is present. Further, referring to FIG. 18 and FIG. 19, the printhead 212 includes plural printheads 218 attached to the carrier substrate 220. In each printhead 218, electrodes 284 for making an electrical connection and an ink supply port 242 are formed on the back surface side of a surface on which a nozzle opening 238 is formed. In the carrier substrate 220 for retaining the printhead 218, electrical wiring is installed on a first surface 270 and second surface 272 thereof, and on the first surface 270 side, the carrier substrate 220 makes the electrical connection with the printhead 218 by solder bumps, and is thus disposed. Further, a logic circuit (not shown) and a drive circuit 230 are mounted on the second surface 272 opposite to the first surface in the substrate 220.\nProblems as described below are inherent in the ink jet recording head formed as described above, in which the ink supply port is formed on the back surface side of the surface of the liquid discharge substrate in which the nozzle opening is formed, the connection electrodes for making the electrical connection with the other members are provided in the vicinity of the ink supply port, and the electrical connection is made with the surface of the carrier substrate on which the electrical wiring is formed.\nFor example, in the ink jet recording head of FIG. 19, the ink supply ports individually formed in the carrier substrate and the liquid discharge substrate must be made to communicate with each other. Specifically, it is necessary to completely prevent the ink from entering the electrical connection portions and completely prevent the ink from leaking to the outside by forming a partition wall for surely separating a fluid on the periphery of the ink supply port.\nFor this purpose, dimensional and positional accuracies of the ink supply port of the liquid discharge substrate and the ink supply port of the carrier substrate become an important point.\nIn the head disclosed in JP 11-192705A, as shown in FIG. 19, the carrier substrate 220 is formed of a plate-like member with some thickness, which is formed of silicon, multilayer ceramics, or glass epoxy, such as one used for forming a hybrid multi-chip module.\nHence, since the head is poor in processability, the positional accuracy of the opening of the ink supply port 242 and the dimensional accuracy of the opening thereof are not very high, and accordingly, a displacement is prone to occur in a relative position thereof to the ink supply port of the liquid discharge substrate. This proves particularly troublesome in the case of using an adhesive or a sealing agent on the fluid partition wall of the ink supply port. In the case of such a construction, a position of the adhesive or the sealing agent is greatly affected by a position of an end of the ink supply port of the carrier substrate, and accordingly, when a position of an end surface of the ink supply port of the carrier substrate is overlapped with the position of the ink supply port of the liquid discharge substrate, there occurs such a problem that the adhesive or the sealing agent flows into the ink supply port of the liquid discharge substrate. Further, in the ink jet head, a necessity to array the plural liquid discharge substrates whose width is as narrow as possible at high density has increased from viewpoints of cost and size. Therefore, it is also necessary to form the ink supply ports of the liquid discharge substrate and the carrier substrate to be narrow in width and at high density. However, in the carrier substrate disclosed in FIG. 19, it is difficult to accurately form the ink supply ports narrow in width because the carrier substrate has some thickness."}
{"text": "The present invention relates to a circuit configuration having a flow controller that can be put into a plurality of states and which outputs a respective command, in a respective one of the states, to a circuit component to be controlled. The flow controller has at least one asynchronously operating delay circuit via which the flow controller moves from one of the states into the respective next state. The invention also relates to an integrated memory and to a test configuration having such a circuit configuration.\nIntegrated memories, for example in the form of DRAMs (Dynamic Random Access Memories), are operated in data processing systems and in this case are actuated by a microprocessor or memory controller, for example. From a certain size of the memory device upward, for example for a memory size of over 1 Mbit, all available DRAMs normally use a so-called xe2x80x9cmultiplex address scheme.xe2x80x9d Such a multiplex address scheme has the advantage that it is very well suited to the functionality of a DRAM.\nIn general, for memory access, first the word lines which are to be addressed and then the bit lines which are to be selected are activated in synchronicity with a clock signal. The address scheme thus involves row addresses being transmitted first and then corresponding column addresses. This selects those memory cells from which data are read and those memory cells to which data are written. Similarly, the microprocessor or memory controller transmits a plurality of single commands, for read access particularly in the form of a word line activation signal, a column access signal and, to complete memory access, a word line deactivation signal.\nOne drawback of such a functionality is, in particular, that a DRAM interface used for the purpose is not immediately aligned with a corresponding processor interface which communicates with fast xe2x80x9cSRAMxe2x80x9d memories. In contrast to DRAMs, SRAMs are very much smaller and faster and they are generally not operated in a multiplex address scheme as described. However, SRAMs have the drawback that they are considerably more expensive than DRAMs.\nEfforts are therefore generally made to use DRAMs also for applications in which SRAM memories have been used previously. In order to be able to use a DRAM also instead of an SRAM, access commands are combined into a macro command which is then processed within one clock cycle. When the macro command has been applied, the single commands for controlling access in a flow controller which can be put into a plurality of states are generated out of sync with the clock signal, with a respective single command being output to the access controller in a respective one of the states. The duration of the individual commands is set using asynchronously operating delay circuits via which the flow controller moves from one of the states into the respective next state. This allows a DRAM to be actuated with a command set that is oriented to fast SRAM memories.\nTo establish the operational functionality or serviceability of an integrated memory, a memory that is to be checked is operated in a test configuration in which a test unit is connected to the memory, and the functionality of the memory is checked. Particularly for the purpose of checking the command sequence for memory access, it is advantageous and appropriate if the single commands run in sync with a clock, since normal test units operate in sync with a clock.\nIt is accordingly an object of the invention to provide a circuit configuration of the above-mentioned type which overcomes the disadvantages of the heretofore-known devices and methods of this general type and in which, in a test mode of the circuit configuration, for example, the commands are output in synchronized fashion.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, a circuit configuration, comprising:\na flow controller enabled to assume a plurality of states and for outputting a respective command, in each of the states, to a circuit component to be controlled;\nthe flow controller having at least one asynchronously operating delay circuit, the flow controller moving from one state into a respectively next state via the delay circuit;\na further signal path connected in parallel with the delay circuit, the further signal path having a clock-controlled multivibrator; and\na switch connected to alternatively operate the delay circuit and the further signal path.\nIn accordance with the invention, the delay circuit in the flow controller mentioned initially has a further signal path connected in parallel with it which has a clock-controlled multivibrator. In the novel circuit, the delay circuit and the further signal path can be operated alternatively using a switch. This allows, for example in a test mode of the circuit configuration, the commands output by the flow controller to be output by the flow controller in synchronized fashion. The invention then affords the advantage that the basically asynchronously operating flow controller, which thereby controls the command sequence, particularly in a normal mode, does not need to be modified. In addition, the additional circuit complexity is kept comparatively low by virtue of the provision of the multivibrator.\nThe invention can be used, in principle, in any type of circuit configuration which has a flow controller which can be put into a plurality of states and which outputs a respective command, in a respective one of the states, to a circuit component which is to be controlled. The invention thus allows the command sequence from a flow controller which operates asynchronously per se to be synchronized with a clock in a second operating state.\nWith the above and other objects in view there is also provided, in accordance with the invention, an integrated memory device that comprises a plurality of memory cells, an access controller for controlling an access to the memory cells, and a circuit configuration as outlined above. The flow controller is connected to output the commands to the access controller for controlling the access to the memory cells.\nIn other words, the invention can be used with particular advantage in an integrated memory in which a flow controller is connected, for the purpose of outputting commands, to an access controller for controlling access to memory cells in the memory. Such a memory is suitable, particularly in a normal mode, for controlling memory access using a command sequence oriented to fast SRAM memories, as described in more detail in the introduction. In addition, however, a memory in accordance with the invention allows, particularly for the purpose of testing the memory using a test unit operating in sync with a clock signal, the command sequence for memory access to be synchronized with the clock signal and hence to be rendered checkable by the test unit.\nIn accordance with an advantageous embodiment of the invention, the clock-controlled multivibrator is in the form of a clock-controlled D-type flip-flop circuit. The delay circuit has an input and an output, with an inverter chain being connected between the input and the output in one embodiment of the delay circuit. The further signal path is connected in parallel with the inverter chain between the input and the output. Alternatively, the delay circuit can have another configuration. It generally outputs the input signal at the output, delayed by a particular time.\nWith the above and other objects in view there is therefore provided, in accordance with the invention, a test configuration that comprises the above-outlined circuit configuration and a test unit connected to the circuit configuration. The further signal path is thereby connected via the switch, and the multivibrator is actuated by a clock signal of the test unit.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in a circuit configuration having a flow controller, integrated memory and test configuration having such a circuit configuration, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims."}
{"text": "Chemokines constitute a family of small cytokines that are produced in inflammation and regulate leukocyte recruitment (Baggiolini M. et al., Adv. Immunol. 55: 97-179 (1994); Springer, T. A., Annu. Rev. Physiol. 57: 827-872 (1995); and Schall, T. J. and K. B. Bacon, Curr. Opin. Immunol. 6: 865-873 (1994)). Chemokines are capable of selectively inducing chemotaxis of the formed elements of the blood (other than red blood cells), including leukocytes such as neutrophils, monocytes, macrophages, eosinophils, basophils, mast cells, and lymphocytes, such as T cells and B cells. In addition to stimulating chemotaxis, other changes can be selectively induced by chemokines in responsive cells, including changes in cell shape, transient rises in the concentration of intracellular free calcium ions ([Ca2+]i), granule exocytosis, integrin upregulation, formation of bioactive lipids (e.g., leukotrienes) and respiratory burst, associated with leukocyte activation. Thus, the chemokines are early triggers of the inflammatory response, causing inflammatory mediator release, chemotaxis and extravasation to sites of infection or inflammation.\nTwo subfamilies of chemokines, designated as CXC and CC chemokines, are distinguished by the arrangement of the first two of four conserved cysteine residues, which are either separated by one amino acid (as in CXC chemokines IL-8, γIP-10, Mig, PF4, ENA-78, GCP-2, GROα, GROβ, GROγ, NAP-2, NAP-4) or are adjacent residues (as in CC chemokines MIP-1α, MIP-1β, RANTES, MCP-1, MCP-2, MCP-3, I-309). Most CXC chemokines attract neutrophil leukocytes. For example, the CXC chemokines interleukin 8 (IL-8), platelet factor 4 (PF4), and neutrophil-activating peptide 2 (NAP-2) are potent chemoattractants and activators of neutrophils. The CXC chemokines designated Mig (monokine induced by gamma interferon) and IP-10 (γIP-10, interferon-gamma inducible 10 kDa protein) are particularly active in inducing chemotaxis of activated peripheral blood lymphocytes. CC chemokines are generally less selective and can attract a variety of leukocyte cell types, including monocytes, eosinophils, basophils, T lymphocytes and natural killer cells. CC chemokines such as human monocyte chemotactic proteins 1-3 (MCP-1, MCP-2 and MCP-3), RANTES (Regulated on Activation, Normal T Expressed and Secreted), and the macrophage inflammatory proteins 1α and 1β (MIP-1α and MIP-1β) have been characterized as chemoattractants and activators of monocytes or lymphocytes, but do not appear to be chemoattractants for neutrophils.\nCC and CXC chemokines act through receptors which belong to a superfamily of seven transmembrane spanning G protein-coupled receptors (Murphy, P. M., Annu. Rev. Immunol., 12: 593-633 (1994); Gerard, C. and N. P. Gerard, Curr. Opin. Immunol., 6: 140-145 (1994)). This family of G-protein coupled (serpentine) receptors comprises a large group of integral membrane proteins, containing seven transmembrane-spanning regions. The receptors are coupled to G proteins, which are heterotrimeric regulatory proteins capable of binding GTP and mediating signal transduction from coupled receptors, for example, by the production of intracellular mediators.\nThe chemokine receptors can be divided into two groups: CC chemokine receptors 1 through 5 (CCR1-5), which bind CC chemokines, and CXC chemokine receptors 1 through 4 (CXCR1-4), which bind CXC chemokines. In general, the CC chemokine receptors occur on several types of leukocytes, and are important for the migration of monocytes, eosinophils, basophils, and T cells (Qin, S., et al., Eur. J. Immunol., 26:640-647 (1996); Carr, M. W., et al., Proc. Natl. Acad. Sci. USA, 91(9):3652-3656 (1994); Taub, D. D., et al., J. Clin. Invest., 95(3):1370-1376 (1995); Neote, K. et al., Cell, 72: 415-425 (1993); Gao, J. -L. et al., J. Exp. Med., 177: 1421-1427 (1993); Charo, I. F. et al., Proc. Natl. Acad. Sci. USA, 91: 2752-2756 (1994); Myers, S. J., et al., J. Biol. Chem., 270: 5786-5792 (1995); Combadiere, C. et al., J. Biol. Chem., 270 (27): 16491-16494 (1995); and Correction, J. Biol. Chem., 270: 30235 (1995); Ponath, P. D. et al., J. Exp. Med., 183: 2437-2448 (1996); and Daugherty, B. L. et al., J. Exp. Med., 183: 2349-2354 (1996); Power, C. A. et al., 1995, J. Biol. Chem., 270: 19495-19500 (1995); Hoogewerf, A. J. et al., Biochem. Biophys. Res. Commun., 218: 337-343 (1996); Samson, M. et al., Biochemistry, 35: 3362-3367 (1996)). In contrast, the two IL-8 receptors, CXCR1 and CXCR2, are largely restricted to neutrophils and are important for the migration of neutrophils (Baggiolini, M., et al., Adv. Immunol., 55:97-179 (1994)). The IL-8 receptors, CXCR1 (IL-8R1, interleukin-8 receptor type 1; Holmes, W. E. et al., Science, 253: 1278-1280 (1991)) and CXCR2 (IL-8R2, interleukin-8 receptor type 2; Murphy, P. M. and H. L. Tiffany, Science, 253: 1280-1283 (1991)) recognize the NH2-terminal Glu-Leu-Arg (ELR) motif, an essential binding epitope observed in CXC chemokines that induce neutrophil chemotaxis (Clark-Lewis, I. et al., J. Biol. Chem., 266: 23128-23134 (1991); Hébert, C. A. et al., J. Biol. Chem., 266: 18989-18994 (1991); and Clark-Lewis, I. et al., Proc. Natl. Acad. Sci. USA, 90: 3574-3577 (1993)).\nIn contrast to monocytes and granulocytes, lymphocyte responses to chemokines are not well understood. Notably, none of the receptors of known specificity appear to be restricted to lymphocytes and the chemokines that recognize these receptors cannot, therefore, account for events such as the selective recruitment of T lymphocytes that is observed in T cell-mediated inflammatory conditions. Moreover, although a number of proteins with significant sequence similarity and similar tissue and leukocyte subpopulation distribution to known chemokine receptors have been identified and cloned, the ligands for these receptors remain undefined. Thus, these proteins are referred to as orphan receptors. The characterization of the ligand(s) of a receptor, is essential to an understanding of the interaction of chemokines with their target cells, the events stimulated by this interaction, including chemotaxis and cellular activation of leukocytes, and the development of therapies based upon modulation of receptor function."}
{"text": "In dialysis apparatuses usually ultrasound-based detectors (hereinafter referred to as detector) are used to protect patients against dangerous air embolisms. For this purpose, a part of an extracorporeal blood circuit to be monitored may be guided in the form of a blood hose through an ultrasonic path of the detectors. For this, a detector includes a piezo element that transmits an ultrasonic (US) pulse or US waves through the blood hose by electric excitation to its resonance frequency. The further detector in the form of a piezo element may convert said US pulse into an electric signal. This signal may be compared to a reference voltage and evaluated by a comparator. On the output side of the comparator a signal is available which informs about whether “air” or “no air” (i.e. fluid or blood only) is provided in the ultrasonic path."}
{"text": "The present invention relates to the attachment of vinyl or metal siding to a building and in particular to a bracket for use in attaching siding to a building.\nMetal and plastic siding are commonly used as an imitation and substitute for wooden lap siding. The metal and plastic siding provide the appearance of wood siding while generally having a longer life span and require less maintenance. Metal siding is typically formed of aluminum or steel while plastic siding is conventionally formed of extruded sections of thermoplastic polyvinyl chloride which is commonly referred to as vinyl siding. In both types of siding, the bottom margin of each panel is typically bent inwardly and then upwardly to form a longitudinal channel with an upstanding inner leg. The top portion of each panel is formed to provide an outwardly and downwardly projecting longitudinal lip corresponding to the channel running along the bottom of the panel positioned thereabove. A securement flange extends above the longitudinal lip of each panel. The panels are typically secured to a wall along their top portions utilizing fasteners driven through the securement flange extending along the top of each panel.\nOne problem common to both metal and vinyl siding is its tendency to expand and contract with changes in temperature. Because of this problem, it has been a practice to incorporate longitudinally extending nail slots along the securement flange through which nails may be driven. It is intended that the nails be driven into the slots a distance sufficient to support the siding but not far enough that the head of the nail would engage the siding pressing it against the side of the building and preventing slidable movement of the siding along the nail through expansion and contraction of the siding. The primary problem with such nail slots is the difficulty in driving a nail or other fastener through the nail slots without fastening the nail too tight thereby preventing slidable movement. This is particularly true when the fasteners are applied utilizing power tools.\nSiding clips as shown in U.S. Pat. Nos. 4,435,933 and 5,150,555 have been developed to overcome the problem of nailing siding too tightly to the wall to which it is secured. Each of the clips generally includes an upper hanger portion which may be nailed to the wall to which the siding is to be secured and a lower portion shaped to receive the upper portion of the siding in a channel or slot while allowing the siding to freely slide laterally through expansion and contraction. One drawback to such clips is that in use it is generally time consuming to slide a large number of clips onto long sections of siding from the ends to provide the appropriate number of clips for supporting the siding. Further, it is then difficult to maneuver the siding with the clips secured thereto prior to securement to the building without having the clips slide off the end of the siding.\nAn additional problem associated with conventional methods for hanging siding which is not corrected through use of siding clips is the tendency of vinyl siding to conform to irregularities or bows of a wall. Because vinyl siding is extremely flexible, the siding will conform to the shape of the wall to which it is attached when it is secured thereto by nailing or through the use of the clips noted above. As a building settles and ages, movement of the foundation and warping, racking or twisting of wood framing members through weathering and natural drying processes results in irregularities and distortions of the wall surface of the building. The materials of construction of a wall to which siding is to be attached may present inherently irregular surfaces such as stucco walls, stone walls or other types of siding. In addition, other factors, such as sloppy construction may cause distortion of the wall surface of a building such that the surface is not flat and is wavy or undulates. Current fastening systems for siding are inadequate to prevent the siding from conforming to such irregularities in the shape of building walls.\nAlthough metal siding tends to be more rigid than vinyl siding, sections of metal siding spanning inward bows on a wall on which it is hung are insufficiently rigid to prevent indentation and permanent deformation of the siding when pressure is exerted against the section of siding spanning such an inward bow. Existing siding attachment systems do not provide sufficient additional rigidity to prevent such permanent deformation.\nAnother drawback of existing siding attachment systems including nailing and the use of clips is uneven sagging. Over time the siding tends to sag. When nails or clips are use to support siding the nails or clips are generally driven into or secured to the framing studs spaced sixteen inches apart such that the siding panels are supported every sixteen inches but not therebetween. Over time the unsupported sections of the siding panel will droop or sag further than the supported sections resulting in an uneven, wavy appearance."}
{"text": "The invention relates to a wheel hub unit comprising a wheel hub, a hydraulic motor and at least one bearing for supporting the wheel hub and the hydraulic motor.\nHydraulic motors for driving wheels of a vehicle are generally known in the art. According to a known wheel hub unit, a steering knuckle is pivotably connected to a kingpin and the hydraulic motor is attached to the steering knuckle at an opposite end of the steering knuckle in relation to the kingpin. The steering knuckle has a cylindrical shape and comprises interior channels for transmitting hydraulic oil to the hydraulic motor. Two bearings are arranged around the steering knuckle at an axial distance from each other.\nA typical range of operation of a hydraulic motor is when the vehicle is running at a speed below a speed limit, for instance below 30 km/h. At higher speeds the hydraulic motor is in a freewheeling mode and a counter pressure is applied for disengaging the pistons in the hydraulic motor. When the vehicle decelerates and the vehicle speed passes said speed limit, the hydraulic motor can be engaged again. At the transition between the freewheeling mode and the operation mode pressure spikes in the millisecond range can occur in which the peak pressure is by far exceeding the normal pressure during operation. Such pressure spikes can reach 100 bar, for instance.\nAccording to the known wheel hub unit, a restriction is provided between the bearings and the bearing which is arranged closest to the hydraulic motor is adapted to let the hydraulic oil pass through to said restriction during said pressure spikes. One problem with this design is that the bearing which is arranged closest to the hydraulic motor will have a relatively short life.\nIt is desirable to provide a wheel hub unit comprising a hydraulic motor, which creates conditions for a long life of the wheel hub unit.\nAccording to an aspect of the present invention, a wheel hub unit comprises a wheel hub, a hydraulic motor and at least one bearing for supporting the wheel hub and the hydraulic motor characterized in that the wheel hub unit comprises a first seal for protecting the at least one bearing from hydraulic oil from the hydraulic motor.\nAdvantageously, by protecting the at least one bearing of the wheel hub unit with the first seal the bearing is not exposed to high pressure lubricant such as hydraulic oil and is particularly not exposed to pressure spikes with high peak pressure. Further, it creates conditions for using a conventional wheel bearing which employs grease as lubricant instead of hydraulic oil from the hydraulic motor. Grease greased bearings can be sealed to the ambient with well-proven seals. Grease greased bearings have favourably a low friction and less abrasive particles are present in the bearing, thus reducing the fatigue of the bearings. Leakage of hydraulic oil through the seal of the bearing or complex seals for sealing the bearing against the ambient can be avoided.\nAccording to a favourable embodiment of the invention, the wheel hub unit comprises means for protecting the first seal from pressure spikes transmittable by the hydraulic oil of the hydraulic motor. This arrangement creates conditions for using soft seals for sealing the bearing section despite high pressure spikes during operation of the hydraulic motor.\nAccording to a favourable embodiment according to the invention, the protection means may comprise a passage between the hydraulic motor and the first seal, wherein hydraulic oil may enter said passage. The passage may be configured for protecting the first seal from pressure spikes transmittable by the hydraulic oil of the hydraulic motor when passing from a freewheeling mode at high speeds to an operation mode at a speed range where the hydraulic motor is operative. The hydraulic oil passage may comprise a region which causes a change of direction of the hydraulic oil for protecting the first seal from said pressure spikes particularly when entering the oil passage in a pressure spike. By changing the direction of the hydraulic oil entering the passage during a pressure spike counteracts the effects of the pressure spike and reduces stress to the components in the oil passage. The change of direction is preferably a relatively sharp change, such as a 90° bend of the passage. The change of direction is preferably at least a 45° bend, and specifically at least a 60° bend of the passage. The passage may be configured with various means, such as for instance, a long extension of the passage, sealing means, restrictions in the passage. Favourably, by protecting the first seal from pressure spikes the first seal can be a standard soft seal. By providing a region causing a change of direction in the hydraulic oil passage a pressure load of the first seal and the effect of pressure spikes can be effectively reduced in the hydraulic oil passage. The hydraulic oil passage may be arranged so that an entrance of the hydraulic oil passage at the hydraulic motor is positioned at a different radial position than an exit of the hydraulic oil passage at the first sealing. Advantageously, the hydraulic oil may enter the hydraulic motor in a central region of the steering knuckle and the exit may be provided at an outer surface of the steering knuckle, resulting in a compact arrangement.\nAccording to a further favourable embodiment according to the invention, the hydraulic oil passage may comprise a first section extending in an axial direction of the wheel hub unit and a second section extending in a radial direction of the wheel hub unit. At the transition between the first and second section the hydraulic oil has to change its direction thus reducing detrimental effects of a pressure spike in the oil passage. The hydraulic oil passage can be designed in an expedient way to reduce pressure spikes.\nAccording to another favourable embodiment according to the invention, the protection means may comprise a second seal which is arranged between the hydraulic motor and the first seal for protecting the first seal from pressure spikes transmittable by the hydraulic oil of the hydraulic motor. Expediently, the second seal may be arranged in said hydraulic oil passage. The second seal can separate the hydraulic oil passage into a high pressure region at one side of the second seal away from the first seal and a low pressure region between the second seal and the first seal. Favourably, the protection of the first seal is improved by the second seal which may be conditioned to seal against a large pressure difference, and the axial length of the hydraulic oil passage can be reduced. This results in a compact arrangement of the wheel hub unit.\nAccording to a favourable embodiment according to the invention, a first drainage channel may be arranged between the first seal and the second seal. Additionally or optionally a second drainage channel may be arranged in front of the second seal. The first drainage channel between the first and the second seal additionally reduces the action of pressure spikes which may pass the second seal. Such oil is discharged to the first drainage channel. The second drainage channel in front of the second seal further reduces an influence of a pressure spike on the first seal. Hydraulic oil entering the hydraulic oil passage is discharged to the second drainage channel thus reducing the pressure load of the second seal. Favourably, the both drainage channels can join in the same drainage channel in the knuckle on which the wheel hub unit is arranged. The first and second drainage channels can be embodied with of without a flow restriction for the hydraulic oil. Expediently, the first drainage channel may be equipped with a restriction device in order to reduce or avoid a pressure load from the joint drainage channel in the knuckle when oil in a pressure spike is discharged to the second drainage channel.\nAccording to a favourable embodiment of the invention, the second seal can comprise a soft seal portion. Expediently, the soft seal portion can provide a defined leakage thus easing requirements for the design of the second seal for a reduction of a pressure spike. By providing a defined leakage, the pressure of the hydraulic oil between the first and the second seal can be controlled. A complex and expensive seal arrangement for providing a fluid tight second seal can be avoided.\nAccording to a favourable embodiment according to the invention, the soft seal portion can be arranged at a radial inner side of a housing which is connected to the wheel hub and configured to rotate with the wheel hub. It is of advantage that the housing can hold the second seal securely in place even during very high pressure spikes. The housing can also take over load caused by the pressure spike thus protecting the second seal from direct actions of the pressure spike. The soft seal portion is configured to rotate with the housing.\nAccording to a favourable embodiment according to the invention, the soft seal portion can be arranged around a spacer ring which can be connected in a rotatably fixed manner to a steering knuckle. The same spacer ring can be used to distribute hydraulic oil into a high pressure section of the hydraulic oil passage. A compact arrangement is achieved.\nAccording to a favourable embodiment according to the invention, the second seal can comprise a leakage channel with a predefined leakage. In an advantageous development, the leakage channel can be arranged between the soft seal portion and the spacer ring. Expediently, a defined pressure level can be established between the second seal and the first seal by a proper design of the leakage channel. As the soft seal portion can rotate with the housing, hydraulic oil in the leakage channel is at the same time a lubricant for the soft seal portion during rotation.\nAccording to a further favourable embodiment according to the invention, a radial inner surface of the first seal is positioned on a radial different, preferably radial larger distance from a wheel hub rotational axis than the leakage channel. The radial different position of the leakage channel with respect to the first seal is of advantage because if pressure spikes transit the channel in unfavourable circumstances, an oil jet may be generated which hits a wall portion of the hydraulic oil channel in the low pressure region of the oil channel and not the first seal. The radial larger distance from the wheel hub rotational axis is of advantage because the first discharge channel is at a radial smaller distance from the wheel hub rotational axis than the leakage channel. As a result, oil can be discharged to the first drainage channel without passing the first seal.\nAccording to a favourable embodiment according to the invention, the wheel hub unit may comprise two bearings and the hydraulic motor may be arranged beside the two bearings in an axial direction of the wheel hub unit. It is of advantage that the two bearings can be constituents of a standard wheel bearing for wheels without a hydraulic motor.\nAccording to a further favourable embodiment according to the invention, the at least one bearing can be formed by a taper rolling bearing. Advantageously, the at least one bearing can be constituents of a standard wheel bearing for wheels without a hydraulic motor.\nAccording to another favourable embodiment according to the invention, the at least one bearing may be configured to be greased by grease. By lubricating the at least one bearing with grease a long lasting bearing arrangement with low friction and low fatigue can be achieved. As such a standard wheel bearing is sealed with soft seals, the invention provides a favourable combination of a hydraulic motor and a standard grease greased wheel bearing.\nAccording to still another favourable embodiment according to the invention, the hydraulic motor may be formed by a radial piston hydraulic motor. A radial piston hydraulic motor has a compact size which is favourable for wheel motors.\nAccording to a further favourable embodiment according to the invention, the first seal may be a soft seal. Favourably, the first seal can be a standard seal as the pressure differences on both sides of the first seal can be maintained comparably small.\nThe invention is particularly useful for commercial vehicles, e.g. tractor vehicles, or construction equipment where wheel motors are helpful for providing a sufficient maneuverability of the vehicle."}
{"text": "Electronic equipment such as mobile phones containing semiconductor devices has been reduced in size and power consumption in recent years. Accordingly, there is a demand for reduction in power consumption in semiconductor devices.\nIt is effective for reducing power consumption to decrease power supply voltages. Decreasing power supply voltages also decreases leakage current, which results in reduction in power consumption. However, decreasing power supply voltages uniformly also decreases available clock frequencies and hence the operating speed. Therefore, multi-supply-voltage semiconductor devices operating at multiple voltages are used in which the power supply voltage supplied to blocks can be independently changed in such a manner that blocks that do not require a fast operating speed from among the blocks constituting the semiconductor device are supplied with a low power supply voltage and blocks that require a fast operating speed are supplied with a high power supply voltage.\nIn a stand-by state, leakage current can be minimized and consequently the power consumption can be reduced by supplying a power supply voltage only to blocks that need to be supplied with the power supply voltage.\nThe term multi-supply-voltage semiconductor device as used in the following description refers to semiconductor devices that operate on multiple power supplies having different voltage values supplied from multiple power supply systems as well as semiconductor devices that operate on multiple power supplies having the same voltage value supplied from multiple power supply systems and semiconductor devices that operate on a single power supply system whose power supply voltage changes.\nIn conventional semiconductor devices consisting of multiple blocks, some or all of the blocks have independent clock circuits, each of which drives a clock signal used within each block in accordance with a clock signal provided from a common clock generator circuit provided for the multiple blocks. However, when a clock signal is driven within each block, clock skew can occur within each block due to timing disagreement in driving the clock signals in the blocks. Such clock skew between blocks can cause a problem of disagreement in signaling timing between blocks. Therefore, semiconductor devices consisting of multiple blocks use a delay circuit to adjust the timing of a common clock signal inputted into the blocks to control clock skew between the blocks.\nFIG. 1 shows a configuration of a conventional multi-supply-voltage semiconductor device that uses such a delay circuit to inhibit clock skew between blocks. The conventional multi-supply-voltage semiconductor device includes two blocks 31, 32, clock (CLK) generator circuit 10, and delay circuit 120. For simplicity, the multi-supply-voltage semiconductor device consisting of two blocks 31, 32 will be described herein. However, semiconductor devices in practice may consist of more than two blocks.\nClock generating circuit 10 generates a clock signal and supplies it to blocks 31 and 32. The clock signal supplied from clock generator circuit 10 to block 32 is delayed by delay circuit 120 inserted between them, and a clock delayed by a certain amount of time from the clock signal generated by clock generator circuit 10 is provided to block 32 as its clock signal.\nBlock 31 includes clock circuit 41 and flip-flop (F/F) circuits 51, 52. Block 32 includes clock circuit 42 and flip-flop circuits 62, 63.\nClock circuit 41 of block 31 drives clock signal CLK1 to be supplied to the circuits within block 31 based on a clock signal from clock generator circuit 10. Clock circuit 42 of block 32 drives clock signal CLK2 to be provided to the circuits within block 32 based on a clock signal delayed by a given amount of time by delay circuit 120.\nIn such a multi-supply-voltage semiconductor device, proper operation must be guaranteed at all operating points even though the device uses a variable power supply which supplies a power supply voltage changing within a certain range as a power supply. That is, the block circuits (clock circuits and arithmetic circuits) must be designed in such a manner that a clock skew is smaller than a signal propagation delay at all power supply voltages, conversely, a signal propagation delay is larger than a clock skew.\nClock skew between blocks 31 and 32 can be suppressed by using delay circuit 120 that provides a certain amount of delay, unless the amount of delays of clock circuits 41, 42 between blocks 31 and 32 changes depending on the voltage value of variable power supply 101.\nFor example, in the conventional multi-supply-voltage semiconductor device shown in FIG. 1, the clock skew between blocks 31 and 32 can be suppressed by setting the amount of delay of delay circuit 120 such that clock signal CLK1 outputted from clock circuit 41 is in phase with clock signal CLK2 outputted from clock circuit 42.\nHowever, if the amounts of delay of clock circuits 41, 42 of blocks 31, 32 do not agree with each other in power supply voltage dependency, a problem arises that a change in the power supply voltage of variable power supply 101 changes each clock circuit delay, increasing the clock skew between the blocks significantly.\nThis problem is more considerable especially if a technique called multi-Vt is used in which MOS transistors with different thresholds (Vt) are formed on the same semiconductor device or if a technique called multi-Tox is used in which MOS transistors having different gate oxide thicknesses (Tox) are formed on the same semiconductor device, because the amounts of delay of clock circuits 41, 42 of blocks 31, 32 differ considerably in power supply voltage dependency from each other.\nFor example, suppose that the power supply voltage dependencies of the amounts of delays of clock circuits 41, 42 shown in FIG. 1 have characteristics as shown in FIG. 2. Even if delay circuit 120 is set so as to prevent clock skew between the blocks at voltage A of variable power supply 101, a delay difference will be produced and therefore clock skew between the blocks will increase when the power supply voltage of variable power supply 101 is changed to voltage B.\nThe timing chart in FIG. 3 shows the operation of clock circuits 41, 42 in such a conventional multi-supply-voltage semiconductor device. Referring to the timing chart in FIG. 3, it can be seen that the clock skew which is minimal at voltage A of variable power supply 101 becomes substantial after the supply voltage of variable power supply 101 is changed to voltage B.\nFurther, as mentioned earlier, some multi-supply-voltage semiconductor devices use both of a non-variable power supply and a variable power supply. FIG. 4 shows a multi-supply-voltage semiconductor device including a block into which a non-variable power supply is inputted and a block into which a variable power supply is inputted. In FIG. 4, non-variable power supply 102 is inputted into block 41 and variable power supply 101 is inputted into block 42. If a different power supply voltage is inputted, a signal with a different level is outputted. In order to accommodate the voltage difference, level shifters 71-73 are provided for signals between blocks 31 and 32. In a multi-supply-voltage semiconductor device as shown in FIG. 4, the power supply voltage of clock circuit 42 in block 32 changes as the supply voltage of variable power supply 101 changes, as shown in FIG. 5. Consequently, clock skew between block 32 and block 31 supplied with different power supplies increases significantly.\nA variety of methods have been proposed for reducing clock skew. Japanese Patent Laid-Open No. 11-39868 discloses, for example, a method for reducing clock skew in a semiconductor integrated circuit system consisting of multiple IC chips. In the semiconductor integrated circuit system, one IC chip is classified as a master chip and the others as slave chips. The master chip detects a change in conditions such as a power supply voltage change and indicates the detected change to each of the slave chips. Each slave chip then adjusts the phase of its clock according to the indicated information about the detected change.\nSince the conventional semiconductor integrated circuit system consists of multiple IC chips, it has a configuration different from that of a multi-supply-voltage semiconductor device in which multiple blocks are formed on a single chip. If the method described above is to be applied to a multi-supply-voltage semiconductor device, a circuit for detecting a change in the power supply voltage and wiring from the circuit to each block are required. In particular, because a multi-supply-voltage semiconductor device has multiple power supply systems, the device requires a circuit for detecting a change in the power supply voltage of each of those power supply systems. Therefore, the number of circuit conductors increases with the number of blocks that are included and power supply voltage systems that are used, which makes it impractical to use this method for high-density semiconductor devices.\nThe conventional multi-supply-voltage semiconductor devices as described above use a delay circuit to control the amount of delay to reduce clock skew at a particular power supply voltage and therefore have the problem that clock skew between blocks increases as the power supply voltage changes if the amounts of delay of each clock circuit provided in each block have a different power supply voltage dependency."}
{"text": "1. Technical Field\nThe present invention relates to a projector.\n2. Related Art\nThere has been known a projector including: a light source device; an optical modulator for modulating a light beam irradiated by the light source device; and a projection optical device for projecting the light beam modulated by the optical modulator in an enlarged manner.\nFor the projector, there has been proposed an arrangement in which an optical diaphragm for adjusting light quantity of the light beam incident on the optical modulator is employed to enhance image quality of a projection image by appropriately changing a brightness of the projection image. The light quantity is adjusted by moving a light shield into and out of an optical path to shield at least a part of the light beam irradiated by a light source. However, it has been difficult to enhance the image quality of the projection image because change over time or manufacturing errors in mechanism components for driving the light shield hinder the light shield from being located at a desired position, in other words, because the projection image cannot be changed to a desired brightness.\nTo solve the problem, there has been proposed a projector employing a light beam detector for detecting the light quantity adjusted by the optical diaphragm, the projector driving the optical diaphragm based on the detected light quantity (see, for example, JP-A-2005-55760). In a case of such projector, even with change over time or manufacturing errors in mechanism components for driving the light shield, the image quality of the projection image can be enhanced, because the light quantity adjusted by the optical diaphragm is detected to reliably detect changes in the light quantity adjusted by the optical diaphragm.\nNoted that, as the light source device for the projectors, light source lamps for generating a light beam by an electric discharge are often used. Such light source lamp for generating a light beam by an electric discharge irradiates a light beam not only in a visible light region but also in a nonvisible light region (such as an infrared region).\nThe projector disclosed in the above-cited document has a problem in using the light source lamp for generating a light beam by an electric discharge, in which, when the light beam including a light beam in the nonvisible light region is irradiated by the light source device on an optical component such as the optical modulator, the temperature of the optical component rises to deteriorate features of the optical component, making it difficult to appropriately maintain the image quality of the projection image."}
{"text": "In our New Zealand Patent Specification No. 229080 (European equivalent Application No. 90304922.9 and Australian equivalent Application No. 55013/90) there is disclosed a process of preparing a hydrolysed lignocellulosic material, the full content of which is hereby incorporated by way of reference.\nThe present invention could have application in such a process/system between different pressure situations or in improved versions thereof. Indeed, the present invention can provide an alternative in many situations to the very expensive and high upkeep \"Bauer valves\" (Sprout-Bauer inlet rotary valves manufactured and sold by Andritz Sprout-Bauer Australia Pty. Limited of 22 Princes Highway, Dandenong, Victoria 3175, Australia), or any such other rotary devices or devices that do not involve rotary chambers and which are used by default in the absence of an economic way of feeding materials from one vessel to another."}
{"text": "Attempts have been made for some time to make the THz frequency range available for imaging applications, in particular in security technology for checking and inspecting people.\nU.S. Pat. No. 5,859,609 discloses a device for holographic approximated real time imaging of concealed articles by means of THz radiation. A vertical linear antenna arrangement is moved over a circular path to achieve cylindrical scanning of 360° of a covered target such as for example a clothed person. The data are present in the form of a non-focused millimeter wave illumination which is capable of reaching a covering such as for example of a clothed person and which is capable of penetrating a covering such as for example clothing. The millimeter wave illumination has to be focused or reconstructed mathematically, that is to say by computation, into recognisable images.\nSuch a device suffers from the disadvantage that for producing imaging of a three-dimensional article by means of a synthetic aperture in two dimensions on the one hand a large number of sources and detectors are required in a first direction and in addition the arrangement of sources and detectors has to be moved completely around the object in order also to be able to generate an image with a synthetic aperture in a second direction perpendicular to the row of sources and detectors. That is linked on the one hand to high costs by virtue of the large number of sources and detectors required, but on the other hand it also involves a comparatively long measurement time as the entire arrangement must be rotated once though 360° in a circle, for example for recording a complete person."}
{"text": "In recent years the automobile has been sophisticated, which entails incremental use of a variety of rotary torque sensors or rotary angle sensors for sensing rotary torque or a rotary angle of a steering shaft in order to control a power steering device or a braking device.\nOne of the foregoing conventional rotary angle and rotary torque sensing devices is described hereinafter with reference to FIG. 5 which is a sectional view of this conventional device.\nIn FIG. 5, conventional device 20 is formed of the following structural elements: first rotator 1, holder 2, magnet 3, second rotator 4, first ferromagnetic body 5, second ferromagnetic body 6, spacer 7, printed circuit board 8, magnetic sensing element 9, controller 11, coupler 12, third rotator 13, first sensor 14, second sensor 15, magnets 16A, 17A, and magnetic sensing elements 16B, 17B.\nFirst rotator 1 is shaped like a cylinder and rotates together with the steering shaft. Holder 2 is shaped like a cylinder. Magnet 3 is shaped like a cylinder where multiple N-poles and S-poles are alternately and adjoiningly arrayed. Magnet 3 adheres to a lower section of an outer wall of holder 2, which adheres to an upper section of an inner wall of first rotator 1.\nSecond rotator 4 shapes like a cylinder. First and second ferromagnetic bodies 5 and 6 are shaped like a ring respectively, and spacer 7 is also shaped like a ring. Second rotator 4 is placed below first rotator 1. First ferromagnetic body 5 is placed above second ferromagnetic body 6, and these two ferromagnetic bodies adhere to an upper section of second rotator 4 such that they confront an outer wall of magnet 3 with a given space therebetween via spacer 7.\nPrinted circuit board 8 has multiple wiring patterns (not shown) on both the faces, and placed beside, substantially in parallel with, first rotator 1 and second rotator 4. Magnetic sensing element 9 includes a Hall element and is disposed between first ferromagnetic body 5 and second ferromagnetic body 6 on a face confronting magnet 3.\nRotary torque sensing section 10 is formed of first and second ferromagnetic bodies 5 and 6, magnet 3 confronting ferromagnetic bodies 5 and 6, and magnetic sensing element 9. Controller 11, which includes electronic components such as a microprocessor, is formed on printed circuit board 8. Controller 11 is coupled to magnetic sensing element 9.\nBetween first rotator 1 and second rotator 4, coupler 12 is provided such that an upper section of coupler 12 adheres to rotator 1 and a lower section thereof adheres to rotator 4. Coupler 12, e.g. torsion bar, is made of steel and shaped like a pole.\nThird rotator 13 includes a spur gear at its underside. First and second sensors 14, 15 include spur gears on their outer walls. Third rotator 13 is rigidly mounted to a lower end of second rotator 4, and first sensor 14 mates with spur gear of third rotator 13 and second sensor 15 mates with spur gear of first sensor 14.\nMagnets 16A and 17A are mounted at the center of first sensor 14 and at the center of second sensor 15 respectively by insert-molding. Printed circuit board 8 is placed beside and generally in parallel with those magnets 16A and 17A. Printed circuit board 8 includes magnetic sensing elements 16B and 17B such as AMR (anisotropic magnetic resistance) at places confronting magnets 16A and 17A respectively.\nMagnet 16A and magnetic sensing element 16B form rotary angle sensor 16, and magnet 17A and magnetic sensing element 17B form rotary angle sensor 17. Magnetic sensing elements 16B and 17B are connected to controller 11. Rotary angle and rotary torque sensing device 20 is thus formed.\nSensing device 20 discussed above has first rotator 1 coupled to the steering shaft and is mounted below the steering wheel of the automobile. Controller 11 is coupled to an electronic circuit of the automobile via connectors and lead-wires (not shown).\nWith the foregoing structure, turning of the steering wheel entails rotation of first rotator 1, and also entails rotation of second rotator 4 which is connected to first rotator via coupler 12. Therefore, third rotator 13 adhering to the lower end of second rotator 4 rotates. This mechanism allows first sensor 14 mated with third rotator 13 and second sensor 15 mated with first sensor 14 to rotate together with third rotator 13.\nThe rotation of first and second sensors 14 and 15 entails the rotation of magnets 16A and 17A, which are mounted at the centers of sensors 14 and 15 respectively. Since magnetisms emitted from magnets 16A and 17A vary due to the rotation, magnetic sensing elements 16B and 17B sense the varying magnetisms, and then a rotary angle sensing signal is supplied to controller 11, where the rotary angle sensing signal repeats fluctuations and draws like a sine-wave, cosine-wave or saw-tooth wave.\nController 11 carries out a given calculation based on the rotary angle sensing signal, supplied from rotary angle sensors 16 and 17, as well as the number of teeth of respective gears. Controller 11 thus senses the rotary angle of third rotator 13, i.e. the rotary angle of the steering shaft, and the result is supplied to the electronic circuit of the automobile.\nThe rotations of the steering shaft and first rotator 1 entail coupler 12 to twist, and then second rotator 4 starts rotating slightly behind first rotator 1. For instance, smaller rotary torque is required for turning the steering wheel when the automobile is driven, so that the delay of second rotator 4 relative to first rotator 1 is small. To the contrary, greater rotary torque is needed when the automobile is halted, so that second rotator 4 starts rotating with a greater delay relative to first rotator 1.\nThe rotations of first and second rotators 1 and 4 entail magnet 3 adhering to those rotators to rotate, and also prompt first and second ferromagnetic bodies 5 and 6 to start rotating after slight delays from magnet 3. Magnetic sensing element 9 senses variation in the magnetism of magnet 3, formed of multiple N-poles and S-poles placed alternately and adjoiningly to each other, via first and second ferromagnetic bodies 5 and 6, and supplies the sensed variation in the magnetism to controller 11.\nAt this time, magnetic sensing element 9 senses weak magnetism when second rotator 4, to which first and second ferromagnetic bodies 5 and 6 adhere, starts rotating after a small delay from first rotator 1, whereas sensing element 9 senses strong magnetism when second rotator 4 starts rotating after a great delay from first rotator 1.\nBased on magnitude of the magnetism sensed via ferromagnetic bodies 5 and 6 by sensing element 9, controller 11 calculates the rotary torque of first rotator 1, i.e. the rotary torque of the steering shaft. The calculated rotary torque is supplied to the electronic circuit of the automobile. Using this calculated rotary torque and other data, the electronic circuit controls the power steering device, the braking device, and other devices of the automobile, where the other data includes the rotary angle of the rotary shaft discussed previously and supplied from controller 11, and various other data supplied from other sensors including a velocity sensor mounted somewhere in the automobile.\nAs discussed above, in response to the rotary angle or the rotary torque supplied from controller 11, conventional rotary angle and rotary torque sensing device 20 controls not only the effectiveness of the braking device in accordance with an angle of the steering wheel turned by a driver but also the operating force for turning the steering wheel.\nPatent Literature of Unexamined Japanese Patent Application Publication No. 2008-82826 is known to the public as one of related art to the present invention.\nConventional sensing device 20 discussed previously includes magnet 3, first and second ferromagnetic bodies 5 and 6, which work for sensing the rotary torque, and those elements are covered with first rotator 1 and second rotator 2. However, when foreign matters such as iron powders attach to those structural elements 3, 5, or 6, magnetic sensing element 9 senses an error in the magnetism emitted from magnet 3 via first and second ferromagnetic bodies 5 and 6. As a result, sensing device 20 fails to sense accurate rotary torque."}
{"text": "1. Field of the Invention\nThe present invention relates to an arrangement for the degassing small high-aspect ratio drilled holes or vias which are present in panels prior to wet chemical processing of the panels. Moreover, more particularly, the present invention is directed to the degassing of vias which are formed in printed circuit boards prior to the copper plating of the vias in order to remove any air or gas bubbles from the vias tending to inhibit the reliable plating thereof. In addition to the foregoing, the invention is also directed to a method for implementing the degassing of the small high-aspect ratio drilled holes which are provided in the panels or printed circuit boards prior to the wet chemical processing of the latter, such as the copper plating of the surface areas of the drilled holes or vias.\nIn the technology directed to the manufacture of electronic packages, especially modules comprising semiconductor devices and printed circuit boards, the modules are frequently provided with considerable numbers of vias or holes, ranging up to 2000 to 3000 in numbers, which are frequently required to have their surfaces copper plated, to produce PTH\"\"s (plated through holes). In attempting to implement the copper plating of such small diameter holes or vias, which normally each possess diameters within a range of about 0.001 to 0.002 inches at high aspect ratios of 6-8:1, difficulties are at times encountered in being able to satisfactorily completely wet or plate the inside of the holes due to the occluding presence of air or gas bubbles. Consequently, in the event that even if a single hole is not adequately copper-plated, then the entire apertured module comprising of the printed circuit board must be discarded, and resultingly leading to considerable economic losses. With respect to the foregoing, it is of extreme importance to be able to provide an adequate wetting of the holes or vias which are present in the printed circuit boards or modules prior to the effectuating of an electroless copper plating pre-cleaning process step, in order to ensure that all of the interior surfaces of the vias or holes are adequately prepared for the electroless copper plating process. This, in effect, would then prevent any potential internal separation phenomena to occur due to a lack of adequate pre-cleaning of the inner planes of the copper material which being applied to the surfaces areas the vias or holes, inasmuch as, if any air or air bubbles were to remain in the holes during the electroless plating process step, the copper plating would fail to fully deposit on the via or hole surfaces, thereby resulting in internal plating separation causing an open circuit and potential failure of the electronic product being equipped with the specific printed circuit board or module.\nBasically, in order to implement the removal of gasses or air bubbles from a liquid bath environment which is employed in the processing of articles, for example, such as modules comprising printed circuit boards incorporating a multiplicity of holes or vias which are to be copper plated in an electroless copper plating process, it is imperative that a suitable degassing step is implemented. This may be preferably carried out through the utilization of an ultrasonic prewetting in a deionized water or other suitable liquid bath preceding cleaning for the electroless plating process, thereby enabling all of the vias or holes to be degassed; in effect, having air removed and the vias or holes filled with liquid; thereby allowing subsequent process cleansing solutions to easily flow into the respective holes or vias in order to facilitate the electroless copper plating process.\nIn order to implement the comprehensive degassing of the liquids; in effect, the removal of any occluded air or gas bubbles, specific physical conditions must be satisfactorily met. Firstly, the liquid in the bath which is utilized in order to initially prewet the modules or panels containing the vias or small holes, and with the bath being constituted of either deionized water or other suitable liquids, must be fully degassed, preferably through the application of ultrasonic energy or any other process in which normal levels of dissolved air entrained in the liquid within the tank and/or entering the tank must be depleted and removed therefrom. Secondly, the apertured modular panels or printed circuit boards within the tank must be exposed to ultrasonic energy in order to promote the displacement and/or the dissolution of air or gas bubbles from the small holes or vias. In essence, failure to employ both of the foregoing conditions would result in the process being ineffective for adequately prewetting all of the small-diameter high-aspect ratio holes or vias, although potentially for small diameter holes or vias with small aspect ratios, agitation in a fully degassed tank (without ultrasonic energy) may be sufficient to dissolve and eliminate all of the air bubbles which are present in the bath or liquid.\n2. Discussion of the Prior Art\nAlthough numerous processes have been developed in various technologies which are directed to the degassing of liquids utilized in the cleaning of articles of the most diverse types, these have not been found to be fully satisfactorily, although complex systems utilize ultrasonic vibrators and energy generators in order to degass the liquids.\nSasaki et al. U.S. Pat. No. 5,810,037 discloses an ultrasonic treatment apparatus for articles immersed in a cleaning liquid which are subjected to a cavitation phenomenon caused by an ultrasonic vibrator. The cleaning liquid is then conveyed to a degassing vessel and thereafter re-circulated to the cleaning vessel subsequent to gas being dissolved in the degassing vessel and the bubbles separated out and maintained in a vacuum.\nPhilips, et al. U.S. Pat. No. 6,210,470 discloses an ultrasonic gas separator wherein gas bubbles in a moving fluid is collected in a chamber arranged adjacent a conduit conveying the fluid. An ultrasonic transducer is adapted to cooperate with a reflector to create a particular ultrasonic wave pattern so as to remove gas bubbles and enable these to be collected exteriorly of the flow of the fluid.\nLong U.S. Pat. No. 6,071,385 discloses a racking fixture for electrochemical processing wherein gases trapped in components are adapted to be removed through suitable ultrasonic applications.\nChung, et al. U.S. Pat. No. 6,042,635 discloses a method for wetting a filter element, and wherein a vibration device will cause pressurized liquid to have entrained air bubbles in a filter element to be separated therein from a liquid and exhausted from the system.\nOhmi et al. U.S. Pat. No. 6,039,814 discloses a cleaning method employing a degassed cleaning liquid with applied ultrasonics, whereby particular ultrasonic frequencies are utilized to remove air bubbles through the intermediary of water vapor from a degassing device containing a liquid.\nKraus Jr. et al. U.S. Pat. No. 5,834,625 discloses an apparatus and a method for debubbling a discrete sample of a liquid through the intermediary of an ultrasonic device which causes air or gas bubbles to be exhausted from essentially high viscous liquids through the formation of ultrasonic waves.\nNicholls, et al. U.S. Pat. No. 5,653,860 discloses a system for the ultrasonic removal of air bubbles from the surface of an electroplated article, and wherein the mechanical vibrations produced by the ultrasonic transducer results in vibrations removing air bubbles and gases from the article surfaces.\nHackett U.S. Pat. No. 5,368,364 discloses a system for removing bubbles from small cavities consisting of blind via holes in a semiconductor wafer by immersing the article and cavities in a liquid medium. The bubbles are enabled to flow freely to enable the processing within the cavities to commence in a suitable degassed liquid.\nShibano, U.S. Pat. Nos. 4,907,611 and 4,865,060 each disclose ultrasonic washing apparatuses whereby air bubbles are removed from a liquid washing specific articles through the intermediary of an ultrasonic vibrator causing the radiation of ultrasonic waves into a washing bath, and enabling the cavitation resulting therefrom to remove soil from soiled articles located in the washing bath.\nJapanese Patent JP 11 94 388 discloses a metal plating apparatus which incorporates a vibration generator to enable air entrapped in a plating liquid to be removed therefrom\nFurthermore, various Japanese publications, including Japanese Patent Publication 10-056261; Japanese Patent JP 73 26 862; Japanese Patent JP 52 1 384; Japanese Patent J 32 48 494; Japanese Patent JP 12 63 300; and Japanese Patent JP 41 54 192, are directed to various aspects of removing entrained or occluded gas or air bubbles from liquids which are employed in the processing of various articles, such as modules or panels including semiconductor devices, printed circuit boards and the like.\nAlthough the foregoing patent publications each describe various types of arrangements, systems and methods for implementing the degassing of liquids, which through suitable applications may be employed for the pre-wetting of articles and any vias or holes formed therein so as to facilitate further wet chemical processing thereof, such as electroless copper plating or the like among other processes, none of these fulfill the above-mentioned conditions in being able to ensure, in a completely satisfactory manner, the pre-treatment of small diameter holes or vias possessing high aspect ratios to provide for the non-problematic subsequent wet chemical processing, such as copper plating. The intent is to produce desired high quality semiconductor devices comprising modules with printed circuit boards each incorporating a large number of small diameter holes or vias. In particular, a review of the current state-of-the technology, particularly as represented by the prior art, indicates that the pre-wetting or cleaning actions are primarily implemented in vessels wherein these contain the liquids in a generally stagnant or static environment, which frequently may be inadequate to fully remove entrapped air or gas bubbles from the modules, and particularly from internally of the small-diameter high-aspect ratio holes or vias which are present in such semiconductor devices or modules in the form of printed circuit board panels.\nIn order to clearly be able to fully utilize the ultrasonic energy, the application of which to the cleaning of articles or components is well known and developed in the technology, the use thereof in a highly reliable manner in order to eliminate entrapped gas or air from vias or holes in a separate manufacturing process is unique, and in order to achieve a high degree of reliability it is necessary to specify the required conditions at which the process becomes most effective.\nAccordingly, the present invention utilizes a novel concept in the form of an in-line degassing arrangement forming an aperture system in which a prewet tank into which the modules or printed circuit boards containing the numerous small high aspect ratio drilled holes or vias are inserted, are degassed through the intermediary of a continuously flowing liquid, such as but not limited to deionized and essentially degassed water rather than the stagnant or static bath tanks which are presently employed in the technology.\nThe continuous water exchange and flow of the liquid employed by the inventive arrangement not only eliminate the air and gas bubbles from the surfaces of the modules and the vias which are formed therein, but also carries off any contaminates and particulates which are freed from the panel and the via surfaces by means of ultrasonic energy or mechanical vibratory action, thus imparting a further cleaning function to the pre-wetting process. This, in effect, will prevent any drying or dewetting of the articles prior to subsequent electroless copper plating, through which a so-called xe2x80x9ctxe2x80x9d could be encountered at the walls of the vias, resulting in a lack of adequate copper plating thickness and presenting a potential reliability problem in the use of electronic packages having the printed circuit boards or modules installed therein.\nThe inventive arrangement and process depends upon a dissolution of bubbles into the gas-depleted liquid flow of the bath, wherein the generated ultrasonic energy field due to the action of ultrasonic transducers serves primarily to promote acoustic microstreaming within the vias, thus shorting the path of diffusion of the air in the liquid by virtue of circulation. The ultrasonic field which is produced also serves to maintain the degassed condition of the bath into which there are immersed the modules, such as the printed circuit board panels incorporate the vias, over the course of time, wherein maintaining the degassed condition of the liquid in the tank is actually of greater significance than a direct application of ultrasonic energy to the panels.\nAlthough so-called panel bumping and panel tilting have been employed in the technology in order to remove air from drilled holes prior to and/or during the plating processes, this maybe somewhat effective for larger-sized holes or vias, but remains essentially ineffective for smaller holes or vias, particularly those possessing high-aspect ratios.\nPursuant to the invention, there is thus provided an in-line, continuous flow pre-wetting process wherein the liquid is degassed prior to entering the pre-wetting tank having the panels immersed in a water bath, such as deionized water, liquid/surfactant, caustic solution or the like, and wherein prior to entering the prewetting tank an ultrasonic degassing unit or units will eliminate air or gas bubbles from the incoming flow of liquid.\nThe degassed liquid is then conducted into a bath tank containing one or more panels, such as printed circuit boards, and a surface spray may also be employed in spraying the upper surface of the bath liquid with the degassed liquid to flush surface impurities.\nAn ultrasonic transmitter array may be provided to face the panel or panels in the tank, and enables the freeing and removal of other contaminants and residual amounts of air from the surfaces and vias hereof, whereas a drain facilitates the continuous flow of the treatment liquid from the tank responsive to an overflow of liquid.\nAccordingly, it is an object of the present invention to provide a novel arrangement for the degassing of small high-aspect ratio drilled holes or vias in panel structures prior to any wet chemical processing of the latter.\nAnother object of the present invention is to provide an arrangement for the degassing of modules comprising panels each incorporating a plurality of small holes or vias through the intermediary of suitably degassed liquid and continuous circulation of the liquid through a bath tank in which the panels are immersed.\nYet another object of the present invention resides in the provision of an arrangement of the type described wherein the degassing of the panels containing the small sized vias or holes prior to wet chemical processing thereof is implemented through the intermediary of ultrasonic energy in a continuous manner inhibiting the de-wetting of the via or holes which have been previously wetted in a capillary-like manner.\nA further object of the present invention is to provide a method for the degassing of small high-aspect ratio drilled holes or vias in panel structures, such as printed circuit boards prior to wet chemical processing thereof."}
{"text": "Memory cells of electronic memory devices such as dynamic random access memory (DRAM) typically employ a number of electronic components in order to store data. For example, in order to store data, a combination of electronic components such as transistors, diodes, and/or capacitors are typically employed in such memory cells to store or not store electrical charges. A charge stored in such a memory cell, may indicate a logic “1”, and no charge stored in such a memory cell, may indicate a logic “0”.\nAn alternative approach for storing data is to use memory cells made of phase change material (PCM). Phase change material may include a material that can be placed into at least two physical states—e.g., a crystalline state and an amorphous state—by increasing or decreasing a temperature of the phase change material. By changing the physical state of the phase change material, certain characteristics of the material, such as electrical resistance, may also change. Such properties may be exploited to form memory cells made of phase change material (referred to herein as “PCM memory cells”).\nAs devices including PCM memory cells (referred to herein as “PCM memory devices”) scale to smaller sizes, fabrication may be limited by higher aspect ratios of device structures and/or limits of conventional lithography and physical vapor deposition (PVD) processes. Techniques and structures that overcome such challenges may allow continued scaling of PCM memory devices to smaller sizes, or may allow increased memory density of PCM memory devices, or combinations thereof."}
{"text": "Gas seals or dry gas seals of the type defined in the introduction are already known from various publications. Examples of these are WO2014/037150 and WO2014/023581.\nDry gas seals of said type generally comprise components composed of ceramic, in particular manufactured from silicon carbide or tungsten carbide. The ceramic components are normally the stationary and/or rotating slide rings and, if appropriate, adjoining elements. A further example from the patent literature is given by DE 39 25 403 C2.\nFeatures of the shaft seal defined in the introduction are already known from DE 202 17 983 U1, DE 20 2007 012050 A1, DE 10 2011 007071 A1, WO 96/04497 A 1.\nIn addition to the primary sealing of the gas seal, in particular dry gas seal, it is required that both the rotating sealing ring and the stationary sealing ring are sealed off with respect to the rotor or stator (housing and possibly connected components) by means of a substantially static seal. Owing to the possible axial vibrations, it is necessary for at least one of these seals to be axially movable. This is generally at least the stationary seal because the mobility is easier to realize at this component. Owing to the high dynamics of the movements, this secondary seal is subjected to extreme loading.\nIt has been found that, in conventional designs, the secondary seal, which generally has to ensure an axial movement, was cumbersome, relatively ineffective and susceptible to defects."}
{"text": "The present invention relates to storage batteries. More specifically, the present invention relates to battery monitors of the type used to monitor individual storage batteries or banks of storage batteries.\nIndividual storage batteries, and banks of storage batteries, are used in various applications including backup power supply applications. For example, remote cellular stations, electrical switching stations, hospitals, and many other installations require a source of backup power. In many such installations, it is important to ensure that the battery or batteries have not degraded and are capable of maintaining a desired amount of charge.\nIn the past, a technician has been typically required to perform battery tests on each of the batteries at the installation. This can be by manually connecting a battery tester to each of the batteries and taking measurements of the batteries. In another example configuration, test connections are run out to each of the batteries and connected to a central location. This allows the testing procedure to be automated.\nThese prior art techniques are cumbersome and are difficult to implement."}
{"text": "The use of wireless communication has become prevalent, especially in the management of assets, particularly those applications associated with inventory management. For example, the use of RFID tags permits the monitoring of the production line and the movement of assets or components through the supply chain.\nTo further illustrate this concept, a manufacturing entity may adhere RFID tags to components as they enter the production facility. These components are then inserted into the production flow, forming sub-assemblies in combination with other components, and finally resulting in a finished product. The use of RFID tags allows the personnel within the manufacturing entity to track the movement of the specific component throughout the manufacturing process. It also allows the entity to be able to identify the specific components that comprise any particular assembly or finished product.\nIn addition, the use of RFID tags has also been advocated within the drug and pharmaceutical industries. In February 2004, the United States Federal and Drug Administration issued a report advocating the use of RFID tags to label and monitor drugs. This is an attempt to provide pedigree and to limit the infiltration of counterfeit prescription drugs into the market and to consumers.\nSince their introduction, RFID tags have been used in many applications, such as to identify and provide information for process control in filter products. U.S. Pat. No. 5,674,381, issued to Den Dekker in 1997, discloses the use of “electronic labels” in conjunction with filtering apparatus and replaceable filter assemblies. Specifically, the patent discloses a filter having an electronic label that has a read/write memory and an associated filtering apparatus that has readout means responsive to the label. The electronic label is adapted to count and store the actual operating hours of the replaceable filter. The filtering apparatus is adapted to allow use or refusal of the filter, based on this real-time number. The patent also discloses that the electronic label can be used to store identification information about the replaceable filter.\nA patent application by Baker et al, published in 2005 as U.S. Patent Application Publication No. U.S. 2005/0205658, discloses a process equipment tracking system. This system includes the use of RFID tags in conjunction with process equipment. The RFID tag is described as capable of storing “at least one trackable event”. These trackable events are enumerated as cleaning dates, and batch process dates. The publication also discloses an RFID reader that is connectable to a PC or an internet, where a process equipment database exists. This database contains multiple trackable events and can supply information useful in determining “a service life of the process equipment based on the accumulated data”. The application includes the use of this type of system with a variety of process equipment, such as valves, pumps, filters, and ultraviolet lamps.\nAnother patent application, filed by Jornitz et al and published in 2004 as U.S. Patent Application Publication No. 2004/0256328, discloses a device and method for monitoring the integrity of filtering installations. This publication describes the use of filters containing an onboard memory chip and communications device, in conjunction with a filter housing. The filter housing acts as a monitoring and integrity tester. That application also discloses a set of steps to be used to insure the integrity of the filtering elements used in multi-round housings. These steps include querying the memory element to verify the type of filter that is being used, its limit data, and its production release data.\nDespite the improvements that have occurred through the use of RFID tags, there are additional areas that have not been satisfactorily addressed. For example, there are a number of applications, such as integrity testing and protein monitoring, in which real time monitoring of the concentration of a particular substance would be extremely beneficial. While RFID tags offer one embodiment of the present invention, solutions utilizing wired communication are also envisioned."}
{"text": "Many electrical connectors enable the conductors of a cable to be electrically joined to terminals in the connector. The connector with the cable extending therefrom can be mated with a compatible connector which is mounted to another cable, an electrical apparatus or the like. Electrical connectors desirably should achieve many objectives including: a high quality electrical connection; low manufacturing costs; a minimum number of components; a low design profile; easy assembly; environmental sealing; compatibility with other connectors and durability.\nOne very desirable electrical connector that achieves the above identified objectives is shown in patent application Ser. No. 005,045 entitled \"ENVIRONMENTALLY SEALED ELECTRICAL CONNECTOR\", now U.S. Pat. No. 4,758,174, which was filed by Leonard H. Michaels and which is assigned to the assignee of the subject invention. The disclosure of this co-pending application is incorporated herein by reference.\nAs explained in application Ser. No. 005,045, it is often necessary or desirable to have the cable extend from the connector at approximately 90.degree. to the mating axis of the connector. However, the mere uncontrolled bending of a cable through 90.degree. immediately adjacent to a connector can create stresses within the cable or at the electrical and mechanical connections between the cable conductors and the terminals of the connector. The stresses created at these locations by the uncontrolled bending of the cable has been known to cause sufficient stress to damage the conductors and render the connector substantially inoperative.\nOne type of structure intended to avoid these problems is a plastic right angle head molded directly onto the assembled connector. Although this prior art structure is compact, it is also costly and suffers from poor reliability. In particular, the insert molding process required for these right angle heads is time consuming and labor intensive. Additionally, these right angle heads employ elastomers as the insulation material to bond to and seal the cable under flexing. However, these elastomers contribute to two other failure modes. For example, the elastomer stresses under tension and becomes proportionally thinner. As a result, it is possible to break the wires under the elastomer at the crimp joint in the right angle head. Damage of this type cannot be seen on the outside of the part, and will only be realized by a failure of the connector in testing or actual use. It also has been found that the bonding of the elastomer to the cable is unpredictable and appears to vary in accordance with unpredictable molding process parameters. In many instances, the elastomer will fail to seal and/or bond to the cable. Consequently, the cable will slip relative to the elastomer, and the intended strain relief will not be provided.\nAnother distinct structure for providing right angle connectors is to employ cast metal angles or multiple cast parts to obtain the required angular orientation. These designs generally work well and do not suffer from the failures described above with respect to the molded on right angle heads. However, connectors employing multiple cast parts are extremely expensive to manufacture and are much larger. Size often is one of the most important design criteria of the connector.\nStill another structure for obtaining the right angle orientation of the connector is to employ multiple metallic parts that are brazed together. Once again, the resulting parts are reliable, but they are costly and labor intensive to manufacture, and are undesirably large.\nThe above described U.S. patent application Ser. No. 005,045 shows and describes a right angle adapter that is molded from an elastomeric material. This adapter is desirable in that it can readily be added on after the cable and connector half are in a fully assembled condition. Although this right angle adapter offers certain advantages, it is desirable to provide a right angle connector that reliably provides more positive strain relief and that securely grips the cable in an unstressed condition.\nIn view of the above, it is an object of the subject invention to provide an electrical connector with a cable exit at an angle to the mating axis of the connector.\nIt is another object of the subject invention to provide an efficient electrical connector with a cable exit at 90.degree. to the mating axis of the connector.\nIt is a further object of the subject invention to provide an electrical connector with a right angle cable exit and that is simple in construction, and easy to assemble.\nIt is still another object of the subject invention to provide an electrical connector with a right angle cable exit wherein the connector provides a seal and strain relief to the cable interface.\nA further object of the subject invention is to provide an efficient electrical connector with a right angle cable exit and with a compact design profile."}
{"text": "Field of the Invention\nThe present invention relates to an information processing equipment and a computer-readable recording medium that generate print data for leaves for leaf pressing processing.\nDescription of the Related Art\nConventionally, leaf pressing print has been known for forming characters and designs made of leaves of gold, silver, and the like, on a sheet. Leaf pressing print is performed by applying adhesive material such as varnish on the sheet, and bonding a metal thin film (leaf) with a point where the adhesive material is applied.\nHere, with reference to FIGS. 1 and 2, a print system including a conventional leaf pressing machine will be described. FIG. 1 is a schematic view illustrating an entire construction of the print system including a conventional leaf pressing machine. FIG. 2 illustrates conventional leaf pressing print (when a leaf object does not stride over a boundary).\nAccording to the conventional leaf pressing print, when print data (submitted data 5) including a gold leaf image and a silver leaf image is submitted from a client terminal 1, the submitted data 5 is divided by a controller 2 into three kinds of print data 6 without leaves excluding a gold or a silver leaf object (also referred to as a “leaf pressed image”), print data 7g for gold leaves including a gold leaf object, and print data 7s for silver leaves including a silver leaf object.\nSubsequently, the print data 6 without leaves is printed by a printer 3, and a user sets printed matter 8 in a leaf pressing machine 4. Since only one kind of leaf pressing can be performed by a general leaf pressing machine 4, a gold leaf roll is set in the leaf pressing machine 4 first, and the leaf pressing machine 4 performs additional print (leaf pressing) on the printed matter 8 based on the print data 7g for gold leaves. Subsequently, a silver color leaf roll is set in the leaf pressing machine 4, and the leaf pressing machine 4 performs the additional print (leaf pressing) based on print data 7s for silver leaves to create printed matter 9 on which the gold leaf and the silver leaf are pressed.\nFurther, in recent years, a leaf pressing machine that can mount two leaf rolls has been developed. As illustrated in FIG. 2, two leaf rolls in different colors are arranged adjacent to each other in a vertical direction with respect to a sheet forwarding direction. A boundary C is a boundary part between a gold leaf region 11g (one example of first region) where the leaf pressing processing is performed with a gold leaf roll 10g (first leaf roll), and a silver leaf region 11s (one example of second region) where the leaf pressing processing is performed with a silver leaf roll 10s (second leaf roll), and corresponds to a center of a leaf image 12 of the submitted data. Such a leaf pressing machine mounting two leaf rolls can press the gold leaf and the silver leaf at a time, if a leaf image 12 includes a gold leaf object 13g and a silver leaf object 13s that are located at a left side and a right side respectively holding the boundary C between the gold leaf object 13g and the silver leaf object 13s as illustrated in FIG. 2.\nPatent Literature 1 discloses constructions of a rotational drive section, a bearing device, a winding axis, and the like, to keep stopping a partial leaf roll that does not perform the leaf pressing, in a device provided with a partial leaf roll sending a partial leaf row obtained by dividing a transfer leaf into one or more."}
{"text": "This invention relates generally to the field of electronic devices, and more particularly to an integrated circuit having a doped porous dielectric and method of manufacturing the same.\nIntegrated circuits typically contain a large number of semiconductor devices. Manufacturers of the integrated circuits often wish to decrease the size of the integrated circuits, which allows more circuitry to be placed within the same physical area of a circuit. To help decrease the size of the integrated circuits, the manufacturers may place the semiconductor devices closer together in the integrated circuits. One problem is that capacitance between the semiconductor devices typically increases as the space between the semiconductor devices decreases. The increased capacitance between the semiconductor devices can interfere with the operation of the semiconductor devices and with the operation of the integrated circuits.\nOne approach to decreasing the capacitance between the semiconductor devices involves decreasing the thickness of the semiconductor devices. For example, the manufacturers may reduce the thickness of gates used in transistors. The capacitance between the semiconductor devices is typically proportional to the cross-sectional area of the semiconductor devices. Because thinner semiconductor devices have less cross-sectional area, the capacitance between the semiconductor devices typically decreases. A problem with this approach is that difficulty may be encountered in maintaining the operation of the semiconductor devices. As the semiconductor devices become thinner, the reduced size of the devices may interfere with the ability of the devices to conduct. Eventually, the reduced thickness of the semiconductor devices may prevent the devices from conducting, and the semiconductor devices in the integrated circuits can fail.\nAnother approach to decreasing the capacitance between the semiconductor devices involves lowering the dielectric constant (K) of the insulating material between the devices. For example, oxide may be used as an insulating material in an integrated circuit, and oxide typically has a dielectric constant of approximately four. The capacitance between the semiconductor devices is typically proportional to the dielectric constant of the insulating material. As a result, lowering the dielectric constant of the insulating material reduces the capacitance between the semiconductor devices. A problem with this approach is that, as manufacturers place the semiconductor devices closer together, the dielectric constant of the insulating material may still be high enough to allow the formation of an appreciable amount of capacitance in the integrated circuit. Also, the insulating material may suffer from contamination, such as by metal ions like sodium, that interferes with the operation of the integrated circuit.\nThe present invention recognizes a need for an improved integrated circuit having a doped porous dielectric and method of manufacturing the same. The present invention reduces or eliminates at least some of the shortcomings of prior systems and methods.\nIn one embodiment of the invention, an integrated circuit includes a semiconductor device. The integrated circuit also includes a contact layer disposed outwardly from the semiconductor device and operable to provide electrical connection to the semiconductor device. In addition, the integrated circuit includes a dielectric layer disposed inwardly from the contact layer and outwardly from the semiconductor device. The dielectric layer comprises an at least substantially porous dielectric material doped with at least one dopant.\nIn a particular embodiment of the invention, the semiconductor device comprises a transistor. Also, the dielectric layer may, for example, include an at least substantially porous dielectric material doped with at least one of phosphorus, fluorine, carbon, and boron.\nIn another embodiment of the invention, a method for forming an integrated circuit having an at least substantially doped porous dielectric includes forming a semiconductor device. The semiconductor device comprises at least a portion of a semiconductor substrate. The method also includes forming a dielectric layer disposed outwardly from the semiconductor substrate and surrounding at least a portion of the semiconductor device. The dielectric layer comprises an at least substantially porous dielectric material doped with at least one dopant. In addition, the method includes forming a contact layer disposed outwardly from the dielectric layer and operable to provide electrical connection to the semiconductor device.\nNumerous technical advantages can be gained through various embodiments of the invention. Various embodiments of the invention may exhibit none, some, or all of the following advantages. For example, in one embodiment of the invention, an integrated circuit is provided that uses a doped, at least substantially porous dielectric material disposed between a semiconductor device and a contact layer. The dielectric material may be doped with any suitable dopant material, such as phosphorus and fluorine. The use of a fluorine dopant decreases the dielectric constant of the dielectric material, which helps to reduce the capacitance between different conductive regions in the integrated circuit. The use of phosphorus reduces the effects of metallic contamination in the dielectric material, which helps to reduce or eliminate interference caused by the contaminant in the integrated circuit.\nAnother technical advantage of at least some embodiments of the invention is that the conductive regions in the semiconductor device retain an appropriate amount of thickness, helping to ensure that the integrated circuit operates properly. The conductive regions in the semiconductor device may retain enough thickness to conduct properly, which helps to ensure that the semiconductor device operates properly. The conductive regions need not be reduced in size to the point where it interferes with the ability of the integrated circuit to function.\nOther technical advantages are readily apparent to one of skill in the art from the attached figures, description, and claims."}
{"text": "Meetings held in large rooms involving two or more individuals can be facilitated using an audio conferencing system. Audio conferencing systems typically include some number of microphones, at least one loudspeaker and a base station which may or may not be linked to a network. The microphones can operate to pick up local acoustic audio signals (speech) and transmit the signals to the base station which generally operates to provide session control and to process the audio signals in a number of ways before sending the signals to a far end system to be played. Among other things, the base station can be configured to track the operational state characteristics of each of the portable devices. Recently, some audio conferencing systems have been designed with wireless or portable microphones, speakers and control devices. Such wireless audio conferencing systems afford much greater flexibility with respect to the conduct of a meeting, permitting individuals to participate in a meeting while moving around the conference room and permitting the optimal placement of the speaker so that it can be heard by everyone in the room.\nFIG. 1 is a diagram showing functional components comprising a commercially available room audio system 100. The system 100 can be comprised of a fixed base station 102 and several different types of wireless or portable devices. In this case, the system 100 has a number of wireless microphones 110, one or more loudspeakers 108, a wireless control device 106, and a charging tray 104. The charging tray 104 has a charging receptacle for each of the portable devices such that when a portable device is placed in the tray an electrical potential will be applied to the portable device and it's battery will start charging."}
{"text": "Many technologies exist that have the ability to digitize different types of input. There are two main touch sensing approaches: active and passive. Passive devices include those that are receptive in a predictable manner to resistive, capacitive, acoustic surface wave, or electro-optical variations due to contact or touch by, for instance, a human finger or a stylus. Active devices include inductive and RF devices.\nThe key downside of active approaches is that an explicit object must be used (e.g., a special pen), which is implemented with electronics (and potentially batteries). This is more complex to manufacture. For example, pens augmented with infrared light emitters on their tips can be used on the commercially available Microsoft Surface.\nTherefore, there is a need to construct a passive input tool without the need for electronics (including e.g., indirectly or wireless powered or EM coupled components including coils) or other active components."}
{"text": "The invention relates to a holder for a drinks container, for example a can, a beaker or a cup.\nHolders of that kind for installation in motor vehicles are known per se in an enormous multiplicity of constructional forms. The known holders have a holding part, which has an insertion opening for inserting the drinks container, and a supporting part, for example of a planar shape, or even bar-shaped, arranged underneath and spaced apart from the holding part, on which supporting part the drinks container inserted into the insertion opening stands. It is known to provide the holders with a guide means so that they can be extended, pivoted out or pulled out into a use position. When the holder is not needed, it can be moved into a retracted position, for example in a door panel or a dashboard of a motor vehicle. The guide means can be in the form of, for example, a sliding guide means, which guides the holder so that it can be pulled out in the manner of a drawer, or a pivotal mounting, by means of which the holder can be pivoted out from the door panel, the dashboard or the like.\nStarting from a holder of the kind described above, the invention is based on the problem of making good use of the installation space required for the holder.\nThe problem is solved according to the invention by the features of claim 1. The holder according to the invention has a receptacle, for example a storage compartment or an ash-tray, which is integrated with the holder. xe2x80x98Integratedxe2x80x99 within the context of the invention is understood to mean that, when the receptacle, holding part and supporting part are in the basic position, the receptacle is arranged in the same installation space as that in which the holding part and supporting part are also located. The receptacle is accordingly located between the holding part and the supporting part or, vice-versa, the holding part and supporting part are located within the receptacle. As a result, the same installation space is used for accommodating both the receptacle and also the holding part and supporting part. The installation space required for the holder according to the invention is therefore not larger than, or only slightly larger than, the installation space required for the holding part and supporting part without the receptacle. The invention consequently has the advantage that the otherwise unused installation space for the holding part and supporting part is additionally used for accommodating the receptacle. The holder according to the invention is preferably so constructed that contents placed inside the receptacle do not impede or prevent a drinks container from being inserted into the insertion opening of the holding part and placed on the supporting part. In an embodiment of the invention, the receptacle is arranged at a level between the holding part having the insertion opening for the drinks container and the supporting part for the drinks container.\nThe holding part and the supporting part of the holder according to the invention are preferably connected to one another; they are especially constructed integrally with one another so that they are always moved together. In principle, a common guide means for the holding part and supporting part is sufficient."}
{"text": "1. Field of the Invention\nThe present invention generally relates to interactive voice technologies and, in particular, to systems and methods for providing users with information in audible form.\n2. Description of the Related Art\nInformation provided to users by computing devices typically is provided in a graphical form. For instance, a personal computer can be used to access a Web site that provides a user with images and/or text via a display device. In this visual context, it is known to combine graphical information provided by a variety of sources and to display the combined graphical information to the user. This enables a user to visually scan the combined graphical information and identify a topic about which the user desires to receive additional information. Once identified, the user can select the topic, such as by actuating a corresponding icon, and be provided with information related to that topic. Typically, the additional information is presented in a maximized format, i.e., at least some of the previously identified topics are no longer displayed to the user and/or the information is presented on a larger portion of the display device.\nIt is also desirable to provide information to users in non-graphical form, i.e., audible form. However, creating an audible analogue to the aforementioned visual context has not been adequately achieved. This is because sources of audible information typically provide grammar fragments that, when combined, cannot be interpreted by voice-rendering applications and/or, if renderable, are unnecessarily impractical. For instance, multiple grammar fragments that are received from multiple sources can be stacked (grouped together) to form a valid VoiceXML document. However, a VoiceXML document with stated grammar fragments is not properly formatted for use by a VoiceXML interpreter that typically is used to render (provide voice) corresponding to VoiceXML documents. By way of further example, if each source of audible information provides a grammar fragment that includes “If you would like to access information on ______, please say ‘_______ ’,” and an interpreter could render the stacked fragments, the user would be provided with redundant audible information, e.g., “If you would like to access information on” is repeated for each source."}
{"text": "The present invention relates to a novel and improved chemical-amplification negative-working photoresist composition suitable for the formation of a photoresist layer on an undercoating of an organic anti-reflection film on the surface of a substrate to give a patterned resist layer of high pattern resolution and having an excellently orthogonal cross sectional profile.\nAlong with the trend in recent years toward higher and higher degree of integration in various semiconductor devices, the photolithographically patterned resist layer on a substrate surface is required to have a pattern resolution of as fine as 250 nm or, as a target in the coming generation, as fine as 200 nm. Needless to say, such an extremely fine pattern resolution of the patterned resist layer cannot be accomplished without an innovative improvement in the performance of the photoresist composition which may be a chemical-amplification negative-working photoresist composition.\nThe above mentioned chemical-amplification negative-working photoresist composition is typically formulated with an acid-curable alkali-soluble resin such as a phenolic resin, a radiation-sensitive acid-generating agent and a crosslinking agent for the resin such as an addition product of urea or melamine and formaldehyde. While it is essential for the formation of a latent image of the pattern that the pattern-wise exposure of the photoresist layer to actinic rays is followed by a post-exposure baking (PEB) treatment, it is known that the line width of the patterned resist layer is influenced by the temperature of the PEB treatment.\nWhen the temperature of the PEB treatment deviates out of a certain range, an appropriate line width of the patterned resist layer as desired can no longer be accomplished and the troubles thereby are more serious as the fineness of patterning is increased. Accordingly, it is eagerly desired to develop a negative-working photoresist composition of which the patterning process can be performed with a greatly increased temperature latitude for the PEB treatment or, namely, the line width of the patterned resist layer is little affected by an increase or decrease of the PEB temperature.\nOn the other hand, it is usual that a photoresist layer of a chemical-amplification negative-working photoresist composition is formed not directly on the surface of a substrate but on an undercoating layer of an anti-reflection film of an organic anti-reflection compound formed on the substrate surface in view of the advantage relative to the pattern resolution which can be as fine as 200 nm or even finer.\nIt is sometimes the case, however, that, when a conventional chemical-amplification negative-working photoresist composition is used in combination with an organic anti-reflection coating compound in the form of a solution such as the most typical commercial products of the DUV Series products (each a product by Brewer Science Co.), an excellently orthogonal cross sectional profile of the patterned resist layer can hardly be obtained and the cross sectional profile is more or less trapezoidal or skirt trailing. Accordingly, it is one of the target problems in the development works for negative-working photoresist compositions to obtain a chemical-amplification negative-working photoresist composition capable of being used in combination with an organic anti-reflection coating compound without being influenced in the pattern resolution and cross sectional profile of the patterned resist layer.\nThe present invention accordingly has an object, in view of the above described problems and disadvantages in the conventional negative-working photoresist composition, to provide a novel and improved negative-working photoresist composition capable of giving a photoresist layer on a substrate surface, which exhibits excellent pattern resolution in a wide range of temperature for the formation of a latent image by pattern-wise exposure of the photoresist layer to actinic rays even when used in combination with an organic anti-reflection film.\nThus, the present invention provides a negative-working photoresist composition which comprises, as a uniform solution in an organic solvent:\n(A) 100 parts by weight of an alkali-soluble resin;\n(B) from 0.5 to 20 parts by weight of an onium salt compound of which the anionic moiety is a fluoroalkyl sulfonate anion as a radiation-sensitive acid-generating agent; and\n(C) from 3 to 50 parts by weight of an ethyleneurea compound represented by the general formula \nxe2x80x83in which R1 and R2 are each a hydroxyl group or an alkoxy group having 1 to 4 carbon atoms and R3 and R4 are each a hydrogen atom, a hydroxyl group or an alkoxy group having 1 to 4 carbon atoms, as a crosslinking agent.\nIn addition to the above defined essential ingredients, the composition can further contain (D) from 0.01 to 1.0 part by weight of an aliphatic amine compound and/or (E) from 0.01 to 1.0 part by weight of a compound selected from the group consisting of carboxylic acids, phosphorus-containing oxoacids and esters of phosphorus-containing oxoacids.\nThe invention also provides a photosensitive patterning material which is an integral layered body comprising (a) a substrate, (b) a layer of an organic anti-reflection composition formed on the surface of the substrate having a thickness of 30 to 300 nm and (c) a photoresist layer formed on the anti-reflection film from the negative-working photoresist composition defined above having a thickness of 200 to 500 nm.\nAs is described above, the essential ingredients in the negative-working photoresist composition of the invention include the alkali-soluble resin as the component (A), an onium salt compound as the component (B) which is an acid-generating agent and a specific ethyleneurea compound as the component (C) which is a crosslinking agent for the component (A).\nThis unique formulation of the photoresist composition has been established as a result of the extensive investigations undertaken by the inventors with an object to develop a photosensitive patterning material of which the total thickness of a negative-working photoresist layer and an organic anti-reflection film does not exceed 800 nm leading to an unexpected discovery that this object can well be accomplished by a negative-working photoresist composition of the above mentioned formulation.\nThe component (A) in the inventive negative-working photoresist composition is an alkali-soluble resinous compound which is not particularly limitative and can be selected from a variety of alkali-soluble resins conventionally used in chemical-amplification photoresist compositions. Examples of alkali-soluble resinous compound particularly preferable from the standpoint of obtaining a patterned resist layer having excellent photosensitivity, pattern resolution and cross sectional profile of the patterned resist layer include copolymeric resins having a weight-average molecular weight of 2000 to 4000 and consisting of 60 to 97% by moles of hydroxystyrene units and 40 to 3% by moles of styrene units, copolymeric resins having a weight-average molecular weight of 2000 to 4000 and consisting of 60 to 97% by moles of hydroxystyrene units and 40 to 3% by moles of styrene units, of which from 5 to 30% of the hydroxyl groups in the hydroxystyrene units are substituted by alkali-insoluble groups and polyhydroxystyrene resins having a weight-average molecular weight of 2000 to 4000, of which from 3 to 40% of the hydroxyl groups in the hydroxystyrene units are substituted by alkali-insoluble groups. More preferably, the alkali-soluble resin as the component (A) is a copolymeric resin having a weight-average molecular weight of 2000 to 4000 and consisting of 60 to 97% by moles of hydroxystyrene units and 40 to 3% by moles of styrene units when excellent orthogonality of the cross sectional profile of the patterned resist layer is essential.\nThe alkali-insoluble group mentioned above is a group which has an effect to decrease the alkali-solubility of a basically alkali-soluble resin when the resin is substituted by such groups. Examples of the alkali-insoluble group suitable for the purpose include tertiary-alkoxycarbonyl groups such as tert-butoxycarbonyl group and tert-amyloxycarbonyl group and lower alkyl groups having 1 to 4 carbon atoms such as methyl group, ethyl group, propyl group, isopropyl group, n-butyl group and isobutyl group, of which the lower alkyl group or, in particular, isopropyl group is preferable in view of obtaining a good patterned resist layer under little influences by the ambient conditions.\nThe acid-generating agent as the component (B) in the inventive photoresist composition is a compound capable of releasing an acid by decomposition under irradiation with actinic rays. While a variety of radiation-sensitive acid-generating compounds are known and used in chemical-amplification photoresist compositions, the component (B) in the inventive photoresist composition is a specific onium salt compound of which the anionic moiety is a fluoroalkyl sulfonate anion. Such an onium salt compound is a known compound as disclosed in Japanese Patent Kokai 54-95686, 62-229942 and 2-120366 and elsewhere.\nIt is taught in Japanese Patent Publication 8-3635 that a preferable acid-generating agent in a negative-working photoresist composition for pattern-wise exposure with a KrF excimer laser beam is tris(2,3-dibromopropyl)isocyanurate because this compound has an advantage of high transparency to the KrF excimer laser beams and high pattern resolution of the patterned resist layer obtained by using this compound as the acid-generating agent.\nThis compound, from which a halogenoacid is generated by irradiation with actinic rays, however, is not suitable for use in a photoresist composition containing an ethyleneurea compound as a crosslinking agent because the effective photosensitivity of the photoresist composition cannot be high enough as to be applicable to an actual production line of LSIs.\nBesides the above named halogenoacid-generating compound, sulfonic acid-generating compounds, such as bis(cyclohexylsulfonyl)diazomethane, are also known and used as an acid-generating agent in chemical-amplification photoresist compositions. These compounds are also not suitable for use in combination with an ethyleneurea compound because a patterned resist layer of high pattern resolution can hardly be obtained with a photoresist composition formulated with these compounds as the acid-generating agent and crosslinking agent.\nIt is the unexpected discovery leading to the present invention that the above described various problems can be overcome by combining a specific onium salt compound as the acid-generating agent and a specific ethyleneurea compound as the crosslinking agent in a negative-working photoresist composition to give a patterned resist layer of excellent properties.\nThe anionic moiety of the onium salt compound as the component (B) is a fluoroalkyl sulfonate anion. The fluoroalkyl group thereof can be a partially fluorinated or fully fluorinated alkyl group. The number of carbon atoms in the fluoroalkyl group is not particularly limitative. It is preferable, however, that the fluoroalkyl group is a perfluoroalkyl group having 1 to 10 carbon atoms because of the general trend that the acid strength of the fluoroalkyl sulfonic acid is increased as the degree of fluorination of the fluoroalkyl group is increased and the number of carbon atoms in the fluoroalkyl group is relatively small not to exceed 10.\nOn the other hand, the cationic moiety as the counterpart of the fluoroalkyl sulfonate anion to form the onium salt compound as the component (B) is not particularly limitative and can be selected from conventional ones. Examples of suitable cations include diphenyliodonium cations and triphenylsulfonium cations optionally substituted by one or more of lower alkyl groups such as methyl, ethyl, propyl, n-butyl and tert-butyl groups, di(lower alkyl)monophenylsulfonium cations, lower alkylcyclohexyl 2-oxocyclohexylsulfonium cations and the like.\nExamples of particularly preferable cations include diphenyliodonium cations represented by the general formula \nin which R5 and R6 are each a hydrogen atom, alkyl group having 1 to 4 carbon atoms or alkoxy group having 1 or 2 carbon atoms, such as diphenyliodonium and bis(4-tert-butylphenyl)iodonium cations, triphenylsulfonium cations represented by the general formula \nin which R7, R8 and R9 are each a hydrogen atom, alkyl group having 1 to 4 carbon atoms or alkoxy group having 1 or 2 carbon atoms, such as triphenylsulfonium, tris(4-methylphenyl)sulfonium and tris(4-methoxyphenyl)sulfonium cations, phenyl dialkylsulfonium cations represented by the general formula \nin which each R10 is an alkyl group having 1 to 4 carbon atoms, such as dimethylphenylsulfonium cations and alkyl cyclohexyl 2-oxocyclohexyl sulfonium cations represented by the general formula \nin which R10 has the same meaning as defined above, such as methyl cyclohexyl 2-oxocyclohexylsulfonium cation.\nThe onium salt compounds particularly preferable as the component (B) in the inventive photoresist composition are those formed from these cations and trifluoromethane sulfonate or nonafluorobutane sulfonate anion or, more preferably, those formed from the triphenylsulfonium cation expressed by the above given general formula (III) and trifluoromethane sulfonate or nonafluorobutane sulfonate anion. These onium salt compounds can be used either singly or as a combination of two kinds or more.\nThe amount of the onium salt compound as the component (B) in the inventive photoresist composition is in the range from 0.5 to 20 parts by weight per 100 parts by weight of the alkali-soluble resin as the component (A). When the amount of the component (B) is too small, the photosensitivity of the photoresist composition cannot be high enough. When the amount of the component (B) is too large, on the other hand, the photoresist composition suffers a decrease in the focusing depth latitude or in the storage stability.\nThe third essential ingredient in the inventive photoresist composition is the component (C) which is an ethyleneurea compound represented by the general formula (I) to serve as a crosslinking agent of the resinous ingredient as the component (A). The compound of the general formula (I) is an ethyleneurea compound substituted at one or both of the N-atoms by a crosslink-forming group or groups selected from hydroxyalkyl groups and lower alkoxyalkyl groups or an ethyleneurea compound substituted at the 4- or 5-position thereof by a hydroxyl group or a lower alkoxy group.\nWhile it is a conventional formulation of negative-working chemical-amplification photoresist compositions that an alkoxymethylated melamine or an alkoxymethylated urea is used as the crosslinking agent, these conventional crosslinking agents have a problem that the crosslinking activity thereof is too high so that, when extremely fine patterning is desired of the photoresist composition, the temperature range for the post-exposure baking treatment is narrowly limited and, when the photoresist layer is formed on an organic anti-reflection film, the cross sectional profile of the patterned resist layer cannot be excellently orthogonal but appearance of skirt trailing or downward narrowing is sometimes unavoidable in the cross sectional profile of the patterned resist layer.\nThe ethyleneurea compound of the general formula (I) used as the crosslinking agent, i.e. component (C) in the inventive photoresist composition, has only a small number of the crosslinkable groups per molecule so that the crosslinking activity of the crosslinking agent is moderately low. When such a crosslinking agent is used in combination with a specific onium salt compound as the acid-generating agent, the above described problems encountered in the use of the conventional crosslinking agent can mostly be overcome. In addition, unexpected advantages can be obtained by the use of the specific ethyleneurea compound as the crosslinking agent that the alkali-solubility of the resist layer in the areas unexposed to the actinic rays is increased so that formation of an isolated resist pattern can be performed with a decreased amount of scum formation in the development treatment.\nThe ethyleneurea compound represented by the general formula (I) can be easily prepared by the condensation reaction of ethyleneurea with formaldehyde to give a condensation product which can optionally be alkoxylated by the reaction with a lower alcohol.\nExamples of the ethyleneurea compounds suitable as the crosslinking agent in the inventive photoresist composition include dihydroxymethyl ethyleneurea, dimethoxymethyl ethyleneurea, diethoxymethyl ethyleneurea, dipropoxymethyl ethyleneurea, dibutoxymethyl ethyleneurea, 1,3-di(methoxymethyl) 4,5-dihydroxy 2-imidazolidinone, 1,3-di(methoxymethyl) 4,5-dimethoxy 2-imidazolidinone and the like. These ethyleneurea compounds are commercially available in the form of a mixture of the monomer, dimer and trimer (MX-280, a product by Sanwa Chemical Co.).\nThe amount of the ethyleneurea compound as the component (C) in the inventive photoresist composition is in the range from 3 to 50 parts by weight or, preferably, from 10 to 20 parts by weight per 100 parts by weight of the component (A). When the amount of the component (C) is too small, the crosslink formation of the resinous ingredient cannot proceed completely resulting in poor properties of the patterned resist layer. When the amount of the component (C) is too large, the photoresist composition suffers a decrease in the storage stability or decrease in the photosensitivity with eventual formation of a particulate matter during storage.\nIn addition to the above described essential ingredients, i.e. components (A), (B) and (C), it is optional that the inventive photoresist composition is admixed with an aliphatic lower-alkyl or -alkanol amine compound as the component (D) and/or a carboxylic acid or an oxoacid of phosphorus as well as an ester thereof as the component (E). These additional ingredients are known and conventionally employed in negative-working chemical-amplification photoresist compositions in the prior art.\nExamples of the above mentioned aliphatic amine compound as the component (D) include tertiary amines such as trimethylamine, triethylamine, tripropylamine, tributylamine, tripentylamine, triethanolamine and tripropanolamine and secondary amines such as dipropylamine, dibutylamine, dipentylamine and dipropanolamine.\nThe amount of the component (D) in the inventive photoresist composition, when added, is in the range from 0.01 to 1.0 part by weight per 100 parts by weight of the component (A).\nExamples of preferable carboxylic acids as the component (E) in the inventive photoresist composition include malonic acid, citric acid, malic acid, succinic acid, benzoic acid and salicylic acid.\nExamples of the oxoacid of phosphorus or a ester thereof as the other class of the component (E) include phosphoric and phosphorous acids and esters thereof such as phosphoric acid, phosphorous acid, di(n-butyl)phosphate and diphenyl phosphate, phosphonic acid and esters thereof such as phosphonic acid, dimethyl phosphonate, di(n-butyl)phosphonate, phenyl phosphonate, diphenyl phosphonate and dibenzyl phosphonate and phosphinic acid and esters thereof such as phosphinic acid and phenyl phosphinate.\nThe amount of the component (E) in the inventive photoresist composition, when added, is in the range from 0.01 to 1.0 part by weight per 100 parts by weight of the component (A).\nIt is optional that the inventive photoresist composition is admixed with the component (D) alone, component (E) alone or both of the components (D) and (E) in combination.\nThe photoresist composition of the present invention is prepared usually in the form of a uniform solution by dissolving the above described essential ingredients and optional ingredients in an organic solvent which is not particularly limitative provided that the solubility of each ingredient therein is high enough. Examples of suitable organic solvents include ketones such as acetone, methyl ethyl ketone, cyclohexanone, methyl isoamyl ketone and 2-heptanone, polyhydric alcohols and derivatives thereof such as ethyleneglycol, ethyleneglycol monoacetate, diethyleneglycol, diethyleneglycol monoacetate, propyleneglycol, propyleneglycol monoacetate, dipropyleneglycol and dipropyleneglycol monoacetate as well as monomethyl, monoethyl, monopropyl, monobutyl and monophenyl ethers thereof, cyclic ethers such as dioxane and esters such as methyl lactate, ethyl lactate, methyl acetate, ethyl acetate, butyl acetate, methyl pyruvate, ethyl pyruvate, methyl methoxypropionate and ethyl ethoxypropionate. These organic solvents can be used either singly or as a mixture of two kinds or more. It is preferable that the organic solvent is a mixture of propyleneglycol monomethyl ether and propyleneglycol monomethyl ether acetate in a mixing proportion of 50:50 to 80:20 by weight in consideration of the good solubility behavior to the components (B) and (C).\nVarious further additives can optionally be added to the inventive photoresist composition including, for example, surface active agents having effectiveness to improve the film-forming properties of the photoresist composition in the formation of a photoresist layer.\nThe present invention further provides a photosensitive patterning material which is an integrally layered body comprising (a) a substrate plate, (b) an organic anti-reflection film formed on the surface of the substrate plate and (c) a photoresist layer formed on the anti-reflection film from the above described negative-working photoresist composition of the invention described above.\nA variety of substrate plates can be used without particular limitations depending on the intended application. Typical examples of applicable substrate materials include semiconductor silicon wafers with or without an undercoating film of SiON, SiN, Si3N4, polycrystalline silicon and TiN and glass plates having a coating film of a metal such as tantalum and chromium.\nThe anti-reflection film formed on one surface of the substrate is limited to a water-insoluble organic anti-reflection film which is formed, different from an inorganic anti-reflection film usually formed by the chemical vapor-phase deposition (CVD) method of an inorganic material, by coating the substrate surface with a coating solution containing an organic ultraviolet absorbing polymeric compound in an organic solvent followed by drying and a heat treatment. Several commercial products of an organic anti-reflection coating solution are available including a product sold under the trade name of DUV-42 (a product by Brewer Science Co.).\nThe organic anti-reflection coating film is formed on the substrate surface in a thickness in the range from 30 to 300 nm. This thickness of the anti-reflection film as well as the thickness of the photoresist layer formed thereon are each an important factor affecting the quality of the patterned resist layer relative to the pattern resolution and orthogonality of the cross sectional profile thereof. In this regard, the photoresist layer formed from the inventive photoresist composition should have a thickness in the range from 200 to 700 nm or, preferably, from 200 to 500 nm or, more preferably, from 200 to 400 nm. When the thickness of the photoresist layer deviates out of the range, satisfactory quality of the patterned resist layer cannot be ensured relative to the dimensional accuracy and orthogonality of the cross sectional profile to negate the advantages to be obtained by the combined use of the specific components (B) and (C).\nIt is further optional that an additional anti-reflection coating film, which must be soluble in water, is formed on the photoresist layer in a thickness of 35 to 45 nm."}
{"text": "Surface patterns such as fingerprints typically are recorded for later associating a unique fingerprint with a corresponding person. Unfortunately, fingerprints and other surface patterns can be distorted depending on how much pressure is applied when making an image of the pattern. Consequently, it can be difficult for a computer or other digital device to identify a person associated with a particular print even though an image of the person's fingerprint is stored in memory for comparison."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory device and a manufacturing method of the semiconductor memory device.\n2. Related Art\nData reading from a NAND flash memory is performed in a page unit. One page includes 4 k-bit to 64 k-bit data, for example. Therefore, in read operation, data of 4,000 to 64,000 memory cells are read simultaneously. Because cell current flows simultaneously to the 4,000 to 64,000 memory cells at this time, potentials of cell source lines sometimes increase (or is lifted). When the potentials of the cell source lines increase, cell current decreases. Consequently, a threshold voltage of the memory cells increases, and a distribution of the threshold voltage of the memory cells is spread. To suppress the spreading of the threshold voltage distribution, the cell source lines connecting between source drivers and sources of the memory cells preferably have as low resistance as possible. For this purpose, a top-layer wiring layer of low resistance in a multilayer wiring structure on a memory cell array can be used as cell source lines and cell well lines (see Japanese Patent Application Laid-open No. 2005-142493 and Japanese Patent Application Laid-open No. 2006-245547).\nIn recent years, to improve performance of memories, the ABL (All Bit Line) sense scheme has come to be used frequently. In the ABL sense scheme, sense amplifiers are provided corresponding to all bit lines and all bit lines are simultaneously read. Therefore, a large amount of current is consumed in the sense amplifiers. To supply power to the sense amplifiers, wide power source lines need to be arranged from a pad to the sense amplifiers. This increases chip size.\nTo further downscale the chip size, there is considered a layout of forming on the top-layer wiring layer of the memory cell array a VSS (ground potential) power source lines to be connected to the sense amplifiers, as well as the cell source lines and the cell well lines (see Kanda et al. “A 120 mm2 16 Gb 4-MLC NAND Flash Memory with 43 nm CMOS Technology” ISSCC 2008/SESSION 23/NON-VOLATILE MEMORY/23.6 p.p. 430, 431 and 625). To decrease resistance of the wiring formed on the top-layer wiring layer in this way, an installation area of each wiring is preferably increased. To increase each wiring area of the cell source lines, the cell well lines, and power source lines as far as possible, each wiring formed on the top-layer wiring layer can cover a whole upper part of the memory cell array. However, when the wiring covers the whole upper part of the memory cell array, hydrogen is not thoroughly supplied to the memory cells in a forming gas/anneal process performed at a final stage of a memory-chip manufacturing process. This becomes a cause of degradation of a data holding characteristic of the memory cells."}
{"text": "This invention relates generally to the separation of meat from bone in a deboning apparatus, but more particularly, pertains to a deboning apparatus that provides the designed support for all of its operating mechanisms in an axially disposed directly driven auger for inducing the separation of bone from meaty material.\nA variety of prior art devices have been available for use in the separation of one component from another particularly in the food processing industry. For example, the U.S. Pat. No. 2,322,791, to DeBack, discloses the use of a cylinder screen surrounding a compression screw, and having an annular space formed between tapers at one end of the screw for use in maintaining back pressure for inducing the expressing of juice from fruits and vegetables deposited into the shown extractor. This type of extractor is pertinent for disclosing the use of an annular valve like means at the end of its screw for regulating the rate of discharge of the fruit pomace from the extractor during its functioning. The Japanese fish industry has long utilized a sorting apparatus for straining fish meat from its vertebrae structure particularly after the whole fish has been filleted thereby leaving residue meat remaining upon the boney structure of the fish. This type of embodiment is disclosed in the Japanese utility model application publication No. 5757/1965, wherein a screw is disposed within a material feed inlet for conveying the deposited fish meat towards a part of the screw that conveys the material towards a strainer roll that is concentrically arranged around the screw so that the meaty component of the fish will be pressured through the roll while the bones are conveyed towards a pressure or squeezing cylinder so as to maintain back pressure upon the material as it is conveyed along by its screw. The strainer disclosed in this publication shows generally the type of arrangement for pressure cylinders that are currently used in the deboning industry, usually with the entire apparatus itself being projecting usually from a bin that holds the bulk of the boning material prior to its being fed into the strainer assembly as herein defined. One of the main problems with this type of prior art apparatus is that the unsupported aspects of its strainer arrangement generates a high degree of vibration during functioning, which has generally limited the operational characteristics of the roll or screw itself at speeds always well below 500 revolutions per minute. Further problems with this type of prior art strainers is that the motor means providing for rotation of the screw has ordinarily not been directly coupled to the shaft of the screw itself, mainly because thought had not been given to the possibility of locating the motor in a direct drive coupling relationship to one end of the screw, generally because of lack of space thereat, and also because these type of strainers, as previously analyzed, were connected directly into the hopper or bin holding the bulk laden material, and wherein, for sanitary reasons, a motor means generally could not be located. Also, prior art deboners operated at much slower speeds which means that their screws could not be directly coupled to the motor. As a result, all of these prior art strainers generally achieve a turning of their screw by an intermediate gear box, and therefore, they are indirectly driven by the motor or drive means. An example of the more contemporary employment of this type of feature in a deboning apparatus currently being marketed is shown in the U.S. Pat. No. 3,739,994, to McFarland. In this particular patent, a compression type of screw is arranged within a meat grinder segment, while the extension of its screw is likewise formed as a compression type screw that is arranged within a perforated conduit. More significantly, it should be noted that the entire apparatus shown therein is of the cantilevered type, and mounts directly to a bin, at one end, and from which the boney material is passed into the grinder segment of the apparatus, with the compression screw achieving its drive indirectly through a gear box as shown. Once again, this particular apparatus of the prior art developes extremely high pressure throughout its extent during its functioning, and since it is unsupported, except at one end, as by being mounted only at one end, it generates enormous vibrations in the event that it is operated at any speeds in excess of 400 rpm. Furthermore, this type apparatus constantly developes enormous compressive pressures entirely along the length of its screw, thereby causing a more accelerated deterioration and frequent replacement of its perforated conduit.\nIt is, therefore, the principal object of this invention to provide a deboning apparatus that is fully supported by an integral base member throughout its length, particularly being bearing mounted at both of its ends so that it may be operated at excessively high revolutions so as to achieve a quicker functioning and greater capacity separation of meat from its boney component during usage of the apparatus.\nAnother object of this invention is to provide a deboning apparatus wherein its auger means arranged within the pressure and perforated housings is directly driven by a motor means without significant vibrations due to the presence of thrust bearings for absorbing the axial forces generated during functioning of the apparatus.\nA related object of this invention is the provision of bearinged pillow blocks integrally upon the base member of a deboning apparatus for fully supporting the operating auger at its approximate ends to thereby eliminate any induced forces or vibration even while the apparatus operates at excessively high speeds.\nAnother significant object of this invention is to provide a deboning apparatus that can achieve two to three times the meat separation capacity output that can be attained from any other meat separating apparatus currently or heretofore available to the trade.\nAnother object is the provision of a deboning apparatus that is formed having its driving motor directly and axially coupled to achieve a rotation of its auger.\nYet another object of this invention is to provide a deboning apparatus that can operate at speeds in excess of 500 rpm and even up to as high as 1800 rpm.\nA further object of this invention is to provide a deboning apparatus incorporating an auger which is separable into a motor coupling, pressure conveying, and valve sections, which are easily separable from each other so as to facilitate either quick assembly, or separation as for cleaning, maintenance, and parts replacement.\nStill another object of this invention is the provision of a deboning apparatus in which its infeed means can be pivotally moved for providing intake of the bone laden material from a variety of directions thereby adding to the versatility and portability of the apparatus for usage.\nA further principal object of this invention is to provide a deboning apparatus in which the significant pressures to achieve meat separation from its bone are built up almost exclusively within its pressure housing section, with said prepressuring being somewhat maintained and gradually dissipating during conveyance of the boney material through the perforated housing of the apparatus.\nYet another feature of this invention is the functioning of a deboning apparatus wherein the pressure release from the prepressurized bone laden material as it traverses by a perforated housing achieves the expressing of the meat from the housing without necessitating any further build up of pressure within the said housing during functioning of the deboning apparatus thereby lessening the wear upon its said perforated housing.\nA related object and benefit achieved from the aforesaid type of function is the reduced abrasion exerted upon the interior of the said perforated housing thereby reducing its incidence for early replacement.\nYet a further object is the provision of a deboning apparatus wherein prepressuring is achieved through the cooperation of a grinder that couples directly to the infeed region of the pressure housing of the deboning apparatus.\nAnother important object of this invention is the provision of a one piece integral solid metal base member for a deboning apparatus that totally supports all of the operating components of the same.\nAnother object of this invention is the provision of an axially arranged discharging means directly beyond the valved section of the deboning apparatus so as to achieve a rapid removal and disposal of the boney components away from the perforated housing.\nAnother object is to achieve a bone abrading separation of fish or meat from its boney components within a perforated housing section that is achieved mainly through a high speed rotation of its contiguous auger.\nThese and other objects will become more apparent to those skilled in the art upon reviewing the summary of this invention, and upon undertaking a study of the description of the preferred embodiment in view of its drawings."}
{"text": "1. The Field of the Invention\nThis invention relates to joining plastics and, more particularly, to novel systems and methods for joining plastics of dissimilar grade.\n2. The Background Art\nIn spin welding two parts together, also known as rotary friction welding, one part is typically held stationary, while the other part is rotated in relative rotation thereagainst. The heat generated by the friction between the two parts causes portions of both to melt. Accordingly, once relative rotation between the parts ceases, the melted material cool sand solidifies, leaving the parts chemically bonded together.\nIt has been discovered that the processes for friction welding or spin welding of polymeric resins of similar grade operate within certain parameters. For example, in conventional spin welding, with proper contact pressure, less than one full turn is required to generate the necessary heat and subsequent melting. These parameters however, are ineffective when resins of differing grades (i.e., cross-grade materials) are involved. The material with the lower melt point or higher melt flow index cannot maintain contact pressure or provide heat sufficient to melt the other material.\nFor example, it has been found that the typical 50 to 200 feet per minute of frictional travel between parts (e.g., 50 to 200 RPM on a 4 inch diameter part) is inadequate to obtain bonding in a cross-grade structural system. Likewise, it has been found that conventional interferences (required to provide friction and the cohesive material) are completely inadequate if the materials in the structure are selected from cross-grade resins.\nAny friction-based welding of cross-grade components results in a preferential melting of the component having the higher melt flow index. That is, the component that melts at a lower temperature, or which flows easier, will preferentially melt, and may completely melt, before any substantial melting has occurred in the other component having a lower melt flow index.\nAccordingly, a new set of operating parameters is required in order to make a friction-based or thermal weld using friction as the thermal mechanism. That is, for example, two materials that are chemically substantially the same, or closely similar in molecular structure may be melted in an injection molding machine or an extrusion machine, and form adequate components or parts. However, friction-based welding, and in particular, spin welding between a molded component (e.g. high melt index) and an extruded component (e.g. low melt flow index) has not been found to work previously."}
{"text": "DC offset cancel in a conventional A/D converter is conducted by inputting two analog signals serving as references, detecting a DC offset error and a gain error of A/D conversion, and canceling the errors in digital processing in a subsequent stage.\nHowever, the above-described conventional technique has a problem that two different reference analog signals are needed. Furthermore, precisions are required of the two different reference analog signals. In addition, there is a problem of an error of the A/D converter involving single-ended to differential conversion, when conducting single-ended to differential conversion on an input signal."}
{"text": "The subject invention relates to interelectrode separators for eliminating zinc dendrite bridging between opposite polarity electrodes in secondary alkaline zinc batteries.\nBattery separators function essentially to electrically isolate the positive and negative plates from each other. They prevent direct contact and suppress \"treeing\" or interelectrode dendrite growth which causes shorting of the respective plates. An ideal separator would isolate the plates without inhibiting electrolyte mobility, and without increasing the battery's internal resistance.\nFor many years, regenerated cellulose was the most commonly used separator for secondary zinc alkaline batteries (e.g., silver-zinc, nickel-zinc). In this regard, fibrous sausage casing and unplasticized cellophane (i.e., 0.025-0.5 mm dry thickness) were wrapped at least one time around the zinc electrode. Often several turns of the separator were used to insure adequate dendrite resistance. These cellulosic separator materials were generally utilized in conjunction with an \"interseparator\" which was a loosely woven or bibulous material wrapped or heat-sealed around the positive electrode to prevent the cellulosic separator from contacting the highly oxidizing environment of the positive electrode. Cellulosics performed well in vented, flooded cells but lacked long-term stability in sealed cells and at temperatures above the ambient (i.e., up to 100.degree. C.). Other problems with cellulosic separators involved swelling and degredation due to oxidation.\nMore recently, newer types of separators with better chemical stability and physical properties have become available. One such class of separators has been designated I/O inorganic/organic) separators and comprises a major inorganic component in a highly stable organic polymer binder. Such separators typically include various inorganic oxides, hydroxides, etc. in a fibrillated polytetrafluoroethylene (hereafter PTFE) binder. These materials are highly flexible and quite durable in the alkaline electrolyte of a secondary zinc battery. The process for making such materials essentially consists of the mixing of the PTFE and inorganic materials under mild shear conditions (e.g., liquid blender) with a high-boiling naphtha lubricant at temperatures above 19.degree. C. to cause fibrillation of the PTFE (see HS-10 bulletin by E. I. Du Pont de Nemours & Co., Plastics Dept., 1968). PTFE levels are normally in the 1-25% by weight range.\nIt is an object of the present invention to provide a new, low cost separator which is strong, flexible and particularly useful in preventing zinc dendrite growth between the electrodes of a secondary alkaline zinc battery. This and other objects and advantages of the present invention will become more readily apparent and understood from the following detailed description and illustration thereof."}
{"text": "1. Field of the Invention\nThe present invention relates to a technique of testing an electric circuit in its operating state. More particularly, it relates to a constitution of an apparatus or system which uses a novel probe suitable for effecting diagnosis and analysis of the operation of an integrated circuit such as an LSI including minute wirings and effects a voltage measurement to thereby be adapted for testing the integrated circuit.\nIn development and manufacture of semiconductor integrated circuits (LSIs), it is essential to test a semiconductor element and, where a malfunction of the element is found, check the cause (i.e., effect a fault analysis). With a recent high degree of integration of LSIs and a multiplication of input/output (I/O) pins, however, it has been difficult to effect an exact design verification and a fault analysis only by measuring signals at the I/O pins by means of an LSI tester or the like. Accordingly, it becomes necessary to effect a voltage measurement of a minute wiring in the LSI, an internal diagnosis and an analysis such as a measurement of operational waveforms. For example, there is known a method of contacting a probe having a pointed end portion directly to a measurement point, amplifying a signal detected by the probe and measuring the signal by means of an oscilloscope or the like. The method is simple and a basic approach. However, it is very difficult to measure a voltage at a much minuter internal wiring or internal electrode than the size of the probe. Also, the precision of measurement is not sufficient and there is a possibility of a secondary trouble such as a short circuit between wirings. Thus, the approach possesses a limitation in itself. In view of this, it is demanded to develop a novel measurement method or approach.\n2. Description of the Related Art\nAs an apparatus suitable for a wiring voltage measurement of a minute pattern, there is known an apparatus using an electron beam or a light beam.\nIn an apparatus using an electron beam, the wiring voltage measurement is carried out by applying the electron beam to a minute wiring (measurement point) in a semiconductor integrated circuit and detecting the amount of secondary electrons emitted from the measurement point. Namely, the measurement using an electron beam utilizes the fact that the amount of secondary electrons to be detected correlates with the voltage at the measurement point.\nOn the other hand, as an apparatus using a light beam, for example, an apparatus utilizing an electro-optic effect is known. In the apparatus, the wiring voltage measurement is carried out by applying the light beam to a certain crystal arranged near the measurement point and detecting the amount of polarization of a light beam transmitted through or reflected from the crystal. Namely, the measurement using a light beam utilizes the fact that, when an external electric field (i.e., wiring voltage) is applied to a crystal, the refractive index of the crystal is changed. The principle of the electro-optic effect is disclosed, for example, in the publication: Valdmanis J. A., Electron. Lett. 23, 1308-1310, 1987. In this connection, the measurement of a change of a short time (e.g., below 1 [ns]) in voltage is carried out by means of a so-called sampling method which effects a voltage measurement by a pulsed beam, not a continuous beam.\nAlso, as another apparatus using a light beam, an apparatus utilizing a photoconductive gate is known. A photoconductive gate has characteristics in which the conductive index is very small (gate OFF state) when light is not applied thereto, and in which the conductive index is increased for a short time (gate ON state) when light is applied thereto. Therefore, it is possible to measure a voltage applied to the gate, i.e., wiring voltage by detecting a current in the photoconductive gate by means of the above characteristics.\nAccording to the above apparatus using an electron beam, it becomes necessary to make the cross section of the electron beam small in accordance with the size of a minute measurement point. In this case, since the number of electrons contained in the beam is decreased and accordingly the number of secondary electrons is also decreased, a problem occurs in that the signal to noise (S/N) ratio is deteriorated. To cope with this, it is effective to increase the number of shots of a pulsed electron beam on a measurement point. In this case, however, another problem occurs in that time required for the measurement is prolonged. Also, the like problem is posed where a time width of an electron beam pulse is shortened for the purpose of an improvement in \"time resolution\" of the measurement. In this connection, considering the transit time effect of second electrons (fluctuation of the measurement timing occurring due to a slow transit speed of second electrons), an upper limit of the time resolution is approximately 5 [ps] at present and thus the realization of a higher time resolution than that value is very difficult in principle.\nOn the other hand, according to the apparatus using a light beam utilizing an electro-optic effect, a very high time resolution exceeding 0.5 [ps] is realized and a voltage resolution corresponding to the S/N ratio is also excellent compared with the above apparatus using an electron beam. However, since the space. resolution is determined depending on the wavelength of light, a problem occurs in that it is very difficult to measure a wiring voltage of a very minute pattern. In this connection, although it is at present possible to effect a voltage measurement with respect to a minute pattern of approximately 1 .mu.m or more, it is very difficult to effect a voltage measurement with respect to a minuter pattern than that value. Namely, the apparatus using an electron beam is excellent in respect of \"space resolution\" although it is insufficient in respect of \"measurement time\" and \"time resolution\", while the apparatus using a light beam is excellent in respect of \"measurement time\" and \"time resolution\" although it is insufficient in respect of \"space resolution\". Namely, both types have advantages and defects contrary to each other.\nAlso, according to the apparatus using a light beam utilizing a photoconductive gate, it is possible to obtain a voltage measurement sensitivity ten times or more compared with the apparatus using an electro-optic effect. However, a few problems are posed. For example, where a photoconductive gate is formed by applying a radiation damage to an SOS (Silicon on Sapphier), there is a disadvantage in that the photoconductive gate must be formed in advance on a LSI chip under test. Although there is a case that a photoconductive gate need not be formed in advance on a LSI chip under test, in this case, a problem occurs in that it is impossible to effect a probing to a minute wiring because the size of an electrode pad to be electrically contacted is relatively large. Also, even if the size of the contact electrode can be made small, it is very difficult to effect a precise probing to the contact electrode with respect to a minute wiring of submicron width."}
{"text": "Turbine engines commonly operate at efficiencies ranging from about 30% to about 40%. The operational efficiency is less than the theoretical maximum because of losses that occur in the flow path. One of the major flow path losses is due to the leakage of hot combustion gases across the tips of the turbine blades. In particular, the leakage occurs across a space between the tips of the rotating turbine blades and the surrounding stationary structure such as the ring segments. This spacing is often referred to as the blade tip clearance.\nBlade tip clearances cannot be eliminated because, during transient conditions such as during engine startup or part load operation, the rotating parts (blades, rotor, and discs) and stationary parts (outer casing, blade rings, and ring segments) thermally expand at different rates. As a result, blade tip clearances can actually decrease during engine startup until steady state operation is achieved at which point the clearances can increase, thereby reducing the efficiency of the engine. Thus, there is a need for controlling blade tip clearances in order to maximize the efficiency of a turbine engine."}
{"text": "This invention pertains to a process wherein silicon wafers are processed in a processing chamber, which is pressurized sequentially with different atmospheres, particularly but not exclusively in oxidation of silicon wafers at high temperatures and high pressures.\nIt is known for silicon wafers to be oxidized by a process known as thermal oxidation and practiced at high temperatures, which typically may range from 700.degree. to 1,000.degree. C., and at high pressures, which typically may range from 5 to 25 atmospheres, in a tubular vessel, which has a removable closure. Typically, the tubular vessel is flooded preliminarily with a purging atmosphere as silicon wafers, which are disposed edgewise and uprightly on a boat, are loaded into the tubular vessel, whereupon the tubular vessel is pressurized to a processing pressure with a series of different processing atmospheres, each processing atmosphere displacing a preceding atmosphere except for a residual portion remaining with such processing atmosphere and diminishing in concentration with time.\nIn a typical case, the purging atmosphere is or comprises nitrogen, which may be regarded as inert under these conditions although there are some indications of its reactivity under these conditions, one of the processing atmospheres is or comprises dry oxygen, which is a slow oxidizer, and another of the processing atmospheres is or comprises steam, which is a fast oxidizer. In the same case, dry oxygen can be introduced initially at a fast rate, so as to promote displacement of the purging atmosphere of nitrogen, and subsequently at a slow rate once the purging atmosphere essentially has been displaced, so as to conserve dry oxygen. Dry oxygen effects slow formation of a thin film of silicon dioxide on the wafers. Next, steam may be introduced at a slow rate, which is limited by physical constraints of available apparatus, so as to effect fast formation of the thin film of silicon dioxide. Also, steam may be followed by dry oxygen, which may be followed by nitrogen, whereupon the tubular vessel may be depressurized.\nOne type of known apparatus, in which the process described above for thermal oxidation of silicon wafers can be practiced, is disclosed in U.S. Pat. No. 4,253,417, wherein the tubular vessel is enclosed within an outer vessel, which is pressurized outside the tubular vessel, so as to equalize the pressures outside and inside the tubular vessel except for a small differential so that leakage of the tubular vessel is leakage out, whereby the tubular vessel may be made of quartz, silicon, or other fragile material. Other types of known apparatus, in which similar processes can be practiced, are disclosed in U.S. Pat. No. 4,018,184 and U.S. Pat. No. 4,167,915.\nIn the process described above for thermal oxidation of silicon wafers, as practiced before this invention, nonuniform oxidation of the wafers has been a problem of serious magnitude. It has been found that the thin films of silicon dioxide formed on the wafers when the wafers are disposed edgewise and uprightly in the oxidizing atmosphere have observable gradients of thickness from thicker portions at upper edges of the wafers to thinner portions at lower edges of the wafers. Also, the problem is exacerbated at higher pressures (20 to 25 atmospheres) among the high pressures (5 to 25 atmospheres) discussed above, and at lower temperatures (700.degree. to 800.degree. C.) among the high temperatures (700.degree. to 1000.degree. C.) discussed above.\nBecause the process is thermally activated, whereupon the process proceeds more rapidly at higher temperatures and less rapidly at lower temperatures, prior efforts to deal with the problem have been directed to improved insulating, baffling, and heating measures, by which it has been found that nonuniformity of the temperature of the oxidizing atmosphere enveloping the wafers can be reduced, controlled, or eliminated, so as to alleviate the problem. The temperature of the oxidizing atmosphere at the upper and lower reaches of the tubular vessel can be measured nonintrusively by thermocouples arrayed on the tubular vessel.\nHowever, it has been found that the problem can be alleviated but cannot be eliminated solely by such measures, as it has been found that the thin films of silicon dioxide formed on the wafers are thicker at upper edges of the wafers and thinner at lower edges of the wafers, even if the temperature of the oxidizing atmosphere enveloping the wafers is uniform at the upper and lower edges of the wafers, and even if the temperature of the oxidizing atmosphere is hotter at lower edges of the wafers and cooler at upper edges of the wafers.\nOther efforts to deal with the problem have been directed to means to promote turbulent flow of the oxidizing atmosphere. However, turbulent flow of the oxidizing atmosphere may be difficult to accomplish effectively, particularly if the oxidizing atmosphere is or comprises steam at high pressures.\nFor a constant volume of an ideal gas, the density of the gas is proportional to the number of molecules of the gas and to the molecular weight of the gas. For constant pressure, volume, and temperature, the number of molecules of the gas is independent of the species of the gas. Thus, the density of the gas is proportional to the molecular weight of the gas at constant pressure, volume, and temperature. Thus, the density of a mixture of two, three, or more gases at constant pressure, volume, and temperature is proportional to the sum of the respective products of the molecular weights and the mole fractions of its components.\nHerein, all molecular weights are stated in round numbers, whenever stated. Herein, the molecular weight of a particular atmosphere being one gas refers to its molecular weight, and the molecular weight of a particular atmosphere being a mixture of gases refers to mean molecular weight of the mixture."}
{"text": "The goal of plant breeding is to combine, in a single variety or hybrid, various desirable traits. For field crops, these traits may include resistance to diseases and insects, resistance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination, stand establishment, growth rate, maturity, plant height and ear height, is important. Traditional plant breeding is an important tool in developing new and improved commercial crops."}
{"text": "1. Technical Field\nAn aspect of the present invention relates to an image reading apparatus, which may be installed in a multifunction device (MFD) having a plurality of functions, such as a scanning function, a printing function, and a copying function.\n2. Related Art\nAn MFD having a plurality of functions, such as a scanning function, a printing function and a copying function, may generally be equipped with an image reading apparatus configured to read an image from a sheet of original document. The image reading apparatus may include, for example, a so-called flatbed-typed image reading apparatus, which reads an image from the original document placed still on a so-called platen (e.g., a transparent plate of glass or resin). For another example, the image reading apparatus may be equipped with an auto-document feeder (ADF), which separates original documents placed on a document placement section automatically and conveys the separated original documents one-by-one, and reads images appearing on the separated original documents while the original documents are sequentially conveyed. Moreover, the flatbed-typed image reading apparatus may be equipped with the ADF.\nThe flatbed-typed image reading apparatus with the ADF may include a document feeder unit and a reader unit. The document feeder unit may be arranged on top of the reader unit, and the document feeder unit may generally be openable (and closable) with respect to the reader unit. Therefore, when a user attempts to have the image read from the original document placed still on the platen, which is arranged in an upper part of the reader unit, in other words, when the user selects not to use the document feeder unit, the user opens the document feeder unit with respect to the reader unit to place the original document on the platen manually. On the other hand, when the user selects to use the document feeder unit, the document feeder unit is not moved but is maintained closed with respect to the reader unit.\nMeanwhile, the above-mentioned reader unit may be configured to form a document placement space, in which the original documents to be read are placed, and a document collecting space, in which the original documents having been read are collected. In order to reduce a space required for installation of the image reading apparatus, and in order to reduce dimensions of the image reading apparatus in depth and width, the document placement space and the document collecting space may be formed in vertically and at least partly overlapping arrangement with respect to each other. In other words, in a plane view, the document placement space and the document collecting space may at least partly overlap each other. While the document placement space and the document collecting space are continuous with each other via a document conveyer path, in order to enable the vertically overlapping arrangement, it is necessary that the document conveyer path includes a reversible path, which is curved upward or downward, and in which the original documents are inverted upside-down.\nWith the reversible path, when so-called regular paper, such as commercially available copier paper, is used as the original document to be read, the regular paper can be flexibly conveyed in the document conveyer path without difficulties. Meanwhile, when thicker paper, such as a postcard and business card, is used as the original document to be read, the thicker paper may be curled in the reversible path along the curvature.\nTherefore, an ADF equipped with a regular conveyer path and a straight-formed different conveyer path may be suggested. The auto-document feeder may convey the regular paper in the regular conveyer path, which is formed to effectively reduce the installation space, and the thicker paper in the different straight conveyer path, so that the deformation of the thicker paper is avoided.\nMore specifically, the ADF may be equipped with the regular conveyer path, including straight paths and an upward curved path continuous with the straight conveyer paths. And in the regular conveyer path, the original document is conveyed from the document placement section arranged in a lower position to the document collecting section arranged in an upper position. At the same time, the ADF may be equipped with the straight-formed different conveyer path, which is branched from the regular conveyer path at an upstream position with respect to the curved path. The straight-formed different conveyer path may be formed to be narrower in width than a predetermined width so that a smaller-sized thicker paper with a width narrower than the predetermined width is guided correctly and allowed to pass there-through. In other words, the different conveyer path may allow the paper with a smaller width, which is narrower than the predetermined width, to pass there-through, while a larger-sized regular paper with a larger width, which is wider than the predetermined width, should be restricted from entering the different conveyer path.\nIn order to restrict the larger-sized wider paper from the different conveyer path and allow the smaller-sized narrower paper in the different conveyer path, the ADF may be provided with a rotatable member, which forms a part of the reversible path and may be rotatable to open and close entrance of the different conveyer path. The rotatable member may be rotated in conjunction with an ejection tray, which may be moved to open when used, and in which the smaller-sized narrower thicker paper should be collected.\nIf, however, the original document having a width wider than the predetermined width is placed on the placement section, and when the different conveyer path being narrower than the predetermined width is open to the wider original document, the wider original document may be drawn into the different conveyer path, of which width is narrower than the width of the wider original document. In this regard, the original document may be damaged or jammed in the different conveyer path."}
{"text": "Link aggregation schemes (e.g., EtherChannel) provide a resilient mechanism for fault tolerance, load balancing and scalable bandwidth by grouping a set of physical network interfaces into a logic link aggregation (LAG). Traffic on the logic link is evenly distributed in steady state and dynamically adjusted in response to system events to avoid traffic polarization and black-holing, which would otherwise eliminate the primary benefits of link aggregation.\nWhile dynamic load distribution is essential to cope with link failure and operational changes, the corresponding hardware reprogramming imposes a major challenge on large scale campus and datacenter switch design. As the number of ports in the switch system increases, both the computational load and the frequency of hardware reprogramming grow exponentially, and hence adversely affect the reliability and system responsiveness of link aggregation schemes."}
{"text": "The present invention relates to a combine harvester. More particularly, it relates to a combine harvester having a cutting and conveyance device for crop, a threshing device for separating grain from the crop, a sieve system, and a blower system.\nDE 39 06 186 A1 makes known a combine harvester of the initially stated type, which comprises a cutting and conveyance device for crop, a threshing unit for separating the grain from the crop, and a sieve system extending in the longitudinal direction of the combine harvester. A blower system is disposed upstream of the sieve system, transversely to the longitudinal direction of the combine harvester. Crop that is picked up by the cutting and conveyance device and conveyed to the combine harvester is directed to the threshing device to separate grain from the crop. The threshing device separates a large portion of the grain from the crop. The remaining crop, which comprises mainly straw and chaff, is transferred to the sieve system in order to separate out the portion of grain remaining in the crop. The blower system disposed upstream of the sieve system in the transverse direction of the combine harvester generates an air volume flow to act upon the sieve system with air, thereby assisting in the separation of heavy grains from the lighter-weight components of the crop, such as chaff and straw. The efficiency of the separation is dependent upon the consistency of the supply of the crop to the sieve system. The consistency of the supply is greatly affected when the combine harvester is in a tilted position since gravity causes crop to accumulate on the sieve system.\nDE 39 06 186 A1 provides a blower system comprising a cross-flow fan which extends in the direction of travel across the entire width of the sieve system. The cross-flow fan is segmented by partition walls in order to attain a number of sectors that corresponds to the number of sieve sections of the sieve system, which are acted upon by air by the cross-flow fan. To balance the uneven crop distribution that occurs when the combine harvester is in a tilted position, the blower speed is changed in accordance with the slant.\nA proven disadvantage of this blower system is that changing the blower speed has the same effect across the entire width of the sieve system. The cleaning output during the harvesting operation on the slope is clearly impaired despite a constant flow across the width of the sieve system at low throughput rates, and therefore, due to the force of gravity, the bulk phase dominates in the crop on the downhill section of the sieve system, while the flight phase dominates on the uphill section of the sieve system."}
{"text": "The present invention relates to air lift pumps, particularly air lift pumps for transporting wet particulate matter from a lower level to an upper level.\nU.S. Pat. Nos. 5,173,194 and 4,126,546 disclose deep bed filter units having a central air lift pump for transporting dirty filter media from a bottom region of the filter to a top region of the filter. Both patents disclose a coaxial pipe arrangement wherein air is forced down a first pipe ('194 down the center pipe, '546 down the outer pipe) and is discharged into the second pipe. The rising air provides the motive force to lift dirty filter medium up the second pipe ('194 up the outer pipe, '546 up the inner pipe). These designs therefore utilize a single media transport pipe and a single point of air injection.\nAs the size of such units is increased, the diameter of the airlift media transport pipe must be increased. Large diameter media transport pipes require an excessive amount of air to provide the proper media flow. Additionally, such large pipes may be more expensive and more difficult to install than smaller diameter pipes.\nA single point of air injection may limit the capacity of such filters. The single point must be located sufficiently close to the transport tube intake to draw the dirty media into the tube. As the flow of water, air, and dirty media travels up the transport tube, friction and turbulence reduce the rate of flow. This flow reduction effectively limits the length of the transport tube, thereby limiting the size of the unit.\nSuch air lift systems spray dirty media and water upwards out of the upper end of the transport tube in an uncontrolled manner. The falling media and water causes splashing when it impacts the top surface of the water contained in the wash chamber. Airborne dirty media and water particles containing impurities may be free to leave the wash chamber contaminating the filtrate. Additionally, the water transported upwards with the dirty media contributes to the volume of dirty water which must be rejected as waste water. Therefore, a significant portion of the influent water (infilt) must be disposed of as waste.\nU.S. Pat. No. 5,277,829, which issued on Jan. 11, 1994 to the assignee of the present invention, discloses a deep bed filter unit in which the infilt flows upwardly through a sand bed such that the filtrate accumulates above the sand bed while the dirty sand is continuously regenerated. Dirty filtration media is collected in the lower region of the filter vessel. An air lift pump system is utilized to transport the dirty media from the lower region to wash means located in the upper region of the filter vessel where impurities are washed out of the filter media. The filter utilizes a single air lift pump having two separate air injection points to provide the motive force necessary to transport the dirty filter media. The first air injection point is located in a horizontal section of the transport pipe and establishes an initial flow of the dirty media. The second air injection point is located in a vertical section of the transport pipe and provides scouring of the dirty media and the necessary motive force to transport the dirty media to the wash means.\nThe use of separate air pumps in a single air pump system requires additional piping, fittings, and diffusers. The additional piping may be comprised of multiple pipe segments or intricately bent pipe to avoid interference with other filter components.\nThe above filters employ a single air lift pump having a single suction point for the dirty media. Such single suction air lift pumps generally do not uniformly remove dirty media from the filter lower region. It is common for dirty media to be more efficiently removed from some areas and less efficiently removed from other areas. In addition, much of the energy imparted by the second air injection point is required for lifting the dirty media to the wash means thereby limiting the amount of media scouring."}
{"text": "In Long-Term Evolution (LTE), time-frequency resources can be divided into 1 millisecond (ms) subframes, where each 1 ms subframe includes two 0.5 ms slots, and each slot (with normal cyclic prefix (CP) durations) includes seven single carrier frequency-division multiplexing (SC-FDMA) symbols in the time domain for uplink (UL) communication and seven orthogonal frequency-division multiplexing (OFDM) symbols in the time domain for downlink (DL) communication. In the frequency domain, resources within a slot are divided into physical resource blocks (PRBs), where each resource block spans 12 contiguous subcarriers.\nIn current LTE systems, resources are usually assigned using a 1 ms minimum transmission time interval (TTI) when data is available for transmission between an eNodeB (eNB) and user equipment (UE). This is referred to as dynamic scheduling. Within each scheduled TTI for UL communication, a UE transmits data over a physical uplink shared channel (PUSCH) in PRB-pairs indicated by an uplink grant to the UE from an eNB that schedules the data transmission. For DL communication, the eNB transmits data over a physical downlink shared channel (PDSCH) in PRB-pairs preceded by a DL assignment from the eNB. The DL assignment information is provided to the UE in a control channel, which is generally referred to as a physical downlink control channel (PDCCH).\nDiscontinuous reception (DRX) functionality is used in LTE to allow UEs to conserve power. Without DRX, a UE is always awake and continuously monitors all subframes of a PDCCH for DL assignments. With DRX, the UE powers down a portion of its circuitry during a DRX sleep mode when there are no expected data packets to be received. The eNB configures DRX with a set of DRX parameters shared with the UE. These DRX parameters can be application dependent such that power and resource savings are maximized. The eNB then schedules DL assignments during periods when the UE is actively monitoring for them.\nWhile DRX operation results in power savings, it comes at the expense of latency. On average, a UE using DRX receives transmissions later as compared to when DRX is not used. This delay can be problematic when the transmissions are associated with a time-sensitive application executing on the UE.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of embodiments of the present disclosure.\nThe apparatus and method components have been represented, where appropriate, by conventional symbols in the drawings, showing only those specific details that are pertinent to understanding the embodiments of the present teachings so as not to obscure the disclosure with details that will be readily apparent to those of ordinary skill in the art having the benefit of the description herein."}
{"text": "This invention relates to impact tools such as hammers. More specifically, this invention relates to an improved hammer having exposed metallic impact surfaces with the remainder of said hammer covered by a unitary resilient encasement.\nHammers of many sizes and shapes are available throughout the prior art, and typically comprise a handle connected to a head having at least one impact surface. The specific construction of the hammer handle and head, and even the striking surface, varies widely according to the desired specialized use of the hammer. That is, hammers for some applications are provided with heads formed from steel or the like to provide hard metallic impact surfaces. See, for example, U.S. Pat. No. 3,341,261. Alternately, for other applications hammers are provided with heads formed from Babbitt or other soft compositions such as lead or rubber to provide soft, non-marring impact surfaces. See, for example, U.S. Pat. No. 2,894,550. Further, some hammers are provided with pellet-filled heads for reducing dangerous recoil upon impact. See, for example, U.S. Pat. Nos. 2,604,914 and 2,737,216.\nAnother specialized hammer construction comprises a skeletal hammer head and handle wholly received within a resilient encasement to provide soft impact surfaces. See, for example, U.S. Pat. Nos. 52,696 and 3,844,321. Such encased hammers are advantageous in that they cannot cause sparking when brought in contact with metallic surfaces, and they do not scratch or otherwise mar surfaces during use. Further, hammers having a resilient exterior coating are not susceptible to undesirable chipping during use and thereby have an improved life span over hammers having heads formed from Babbitt or the like.\nResilient encasements have not been used, however, with hammers requiring steel or other hard metal impact surface since it is necessary to keep the hard impact surface exposed. As a result, such hammers are therefore left with extensive exposed steel surface areas over the hammer head, and often over the handle too. These exposed steel surface areas can chip, create sparks, or scratch surface finishes both during use and when the hammer is put down. That is, when the hammer is used near other machinery, etc., portions of the hammer other than the impact surface can contact the machinery to cause sparks, chips, and scratches. Similarly, when the hammer is put down, the portions of the hammer other than the impact surface can contact other metal surfaces to cause sparks, etc.\nIt is desirable, therefore, to provide an improved hammer having exposed steel or other hard-faced impact surfaces wherein the remainder of the hammer is covered by a unitary, resilient encasement. Moreover, it is desirable to provide such an improved hammer which is quickly and easily assembled, made from economical materials, and of durable construction."}
{"text": "Websites that provide electronic services receive input from client machines in the form of requests to access content at various addresses. These requests may follow a particular sequence. For example, within a user session, an online banking service may receive, in order, a request to access, in order, a login page, a password page, an account summary page, a transaction history page, a funds transfer page, a confirmation page, and a logout request page.\nHowever, the online banking service and websites like it are susceptible to attacks from malicious actors. For example, in some attacks, malicious code in a browser running on a user computer alters the requests to a web server in a way that avoids detection by the user.\nConventional approaches to detecting attacks from malicious actors involve comparing frequencies of various behaviors in user sessions with predicted frequencies of those behaviors. For example, an administrator may analyze the frequency of a particular viewing order of a web page and compare this frequency to that predicted by a statistical model. If the frequency in the user session exceeds the predicted frequency, an administrator labels the user session as suspect."}
{"text": "Conventional precision screw assemblies mandate accurate alignment of the nut and screw axes throughout the range of movement. The tables or other objects driven by the screw assembly must also be mounted so as to be accurately guided in their axial or linear motions. Also, deviations in the axis of the screw must be minute."}
{"text": "1. Field of the Invention\nThe present invention relates to a cigarette lighter device mounted on a vehicular floor console or the like.\n2. Description of Related Art\nConventionally, for driver's convenience, a vehicular cigarette lighter has been assembled to an instrument panel, a floor console, or the like. Also, cigarette lighters have been disposed at a side position or an upper position of an ashtray in view of their properties (see Japanese Utility Model Provisional Publication No. 62-6147). As a conventional example, there is also available an ashtray in which a cigarette lighter is incorporated.\nHowever, from a structural viewpoint, the above-described conventional cigarette lighter has a possibility that cigarette ash adhering to the lighter may drop, so that a floor etc. in a cabin may be made dirty. Also, some ashtrays incorporating a cigarette lighter have a mechanism such that cigarette ash adhering to the lighter drops into the ashtray. In this case, the layout must be such that the cigarette lighter and the ashtray are necessarily located at fixed positions, which poses a problem of restricted degree of freedom of design."}
{"text": "1. Field of the Invention\nThis invention relates to the field of material deposition, and more specifically, to the deposition of one or more semiconductor layers/films onto a substrate member.\n2. Description of the Related Art\nIt is known in the field of semiconductor manufacturing that semiconductor material be selectively deposited onto a substrate member. In so doing, it is known that an in-line horizontal, and in-line vertical, or a circular assembly of vacuum deposition chambers may be used wherein the substrate member is sequentially moved from one chamber to another as is needed to achieve a desired sequence of semiconductor layers on the substrate member.\nIt is usual that the limited internal volume of a vacuum deposition chamber limits the physical size of a substrate member that can be inserted into the chamber(s), and upon which semiconductor material may be deposited.\nIt would be desirable to provide a means whereby a substrate holder would enable a larger area of substrate to be operated upon by a deposition chamber(s)."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of extracting characteristic image data, and more specifically, to a method of extracting characteristic image data, such as density data or the like of a human face, which is used when a color original image is copied onto a color copy material or a black-and-white copy material.\nThe present invention also relates to a color data conversion device for an image processing apparatus, and in particular, to a color data conversion device for an image processing apparatus in which color data obtained from a color original image is converted into hue values, saturation values, and lightness values.\n2. Description of the Related Art\nWhen a photograph of a person is viewed, the area which is most noticed is the person's face. In order to produce high-quality photographs, it is necessary to print the color and the density of a human face at an appropriate color and an appropriate density.\nConventionally, a face region in an original image of a color film is designated by a light pen, and the density data of the human face is extracted. The amount of exposure is determined based on the extracted density data so that the color of the face is printed appropriately. Such technology is disclosed in Japanese Patent Application Laid-Open Nos. 62-115430, 62-115431, 62-115432, 62-189456, 62-189457, 63-138340, and 63-178222.\nHowever, in the above-described conventional art, there is a drawback in that the printing operation requires much time because an operator must use the light pen to designate the face region of each image. Further, automatization of the technology is difficult because an operator must view the image and designate the face region.\nMethods of automatically extracting human face data have been proposed. For example, Japanese Patent Application Laid-Open Nos. 52-156624, 52-156625, 53-12330, 53-145620, 53-145621, and 53-145622 disclose the following method in which human face data is extracted by the extraction of flesh color data. Namely, a color original image is divided into a plurality of points to be photometrically measured. Each of the points to be photometrically measured is divided into three colors of R (red), G (green), and B (blue) which are photometrically measured. A determination is made as to whether the color of each photometrically measured point calculated from the photometric data falls within a flesh color range. A cluster (group) of photometrically measured points which have been determined to be within a flesh color range are used for face density data.\nHowever, in this method, because a color within a flesh color range is used as face density data, regions, which are not face regions and which are flesh color regions such as the ground, tree trunks, clothes, and the like or are regions of a color which is similar to a flesh color, are extracted as density data. Further, even if the same subject is photographed under the same conditions, the hue of the photographed image differs in accordance with the type of film used. Therefore, when the type of film differs, there are cases in which the face density data cannot be extracted automatically. Moreover, when the color of the light source which illuminates the subject is different, the hue of the photographed image differs (e.g., an image photographed with fluorescent light as the light source may be greenish). Therefore, when the light source color differs, there are cases in which the face density data cannot be extracted automatically.\nIn order to solve the above-mentioned drawback which arises when the light source color differs, the photometric data of the flesh color ranges may be extracted after light source color correction is effected. Light sources can be divided broadly into sunlight, fluorescent light, and tungsten light. However, when sunlight is the light source, the hue differs in accordance with the season and the time of day. Even if the season and the time are the same, the hue differs in accordance with direct or indirect light. Further, with artificial light sources such as fluorescent light, hues vary greatly as there are various types of manufactured light sources. Accordingly, it is difficult to designate the type of light source and effect light source correction for each light source. Even if color correction could be effected perfectly, extraction could not be effected such that flesh color regions such as the ground, tree trunks or the like or regions of a color similar to a flesh color were not extracted, and cases in which the types of films differ could not be addressed.\nConventionally, there exist image processing apparatuses such as image displaying devices, image recording devices and the like. In an image displaying device, an image is read from a color original image such as a film, a print or the like, and the color original image is reproduced and displayed on the basis of the read data. In an image recording device, on the basis of the read data, an amount of exposure is determined and an image is recorded so that the color original image is reproduced.\nFor example, the density data of a main image, such as a human face or the like, on an original image of a color film is extracted. Based on the extracted density data, an amount of exposure is determined such that the face color is printed appropriately. Such technology is disclosed in Japanese Patent Application Laid-Open Nos. 62-115430, 62-115431, 62-115432, 62-189456, 62-189457, 63-138340 and 63-178222.\nGenerally, in image processing apparatuses such as those described above, in order to facilitate calculation for determining a composite color when color correction of the color original image is effected and in order to facilitate calculation for determining the tone of the color original image, R data, G data, B data expressing the densities of the three colors of red light, green light and blue light (i.e., the three primary colors R, G, B) are converted into H data, S data, and L data respectively expressing hue (H), saturation (S), and lightness (L) which are used in a color development system.\nKnown methods of conversion include conversion by use of a HSL hexagonal pyramid, and conversion based on the definitions of Hayden and Reins.\nHowever, due to complex calculation using various arithmetic expressions, the processing times of the above-mentioned conversions are long. Accordingly, in an image processing apparatus, the processing time required for conversion alone is long. A drawback arises in that when this processing time is added to the time for fundamental processing of the image, the entire processing time of the processing apparatus greatly increases.\nIn order to solve this drawback, the calculating time can be reduced by using a high-speed type CPU. However, this results in increasing the cost of the processing apparatus. Further, the results of calculation for the above-mentioned conversions are stored in tables in advance, and the respective HSL data corresponding to the RGB data can be output by using these tables. In this way, because calculation at the time of conversion is not used, the calculating time is shortened, and processing at the apparatus can be expedited.\nHowever, in order to output the respective HSL data corresponding to the combinations of the respective RGB data, it is necessary to provide respective tables for H, S, L. Further, the number of indices which designate the output data of the tables is the number of combinations of respective R, G, B data. Therefore, the storage capacity is large. Namely, in order to decrease the difference between densities of an image, the gradations of the image are increased. For example, when the gradations of the image are expressed by 8 bits (0 to 255), indices of (2.sup.8).sup.3 bits are required for the respective tables for outputting the HSL data. Accordingly, when the respective HSL output data are 8 bits, the storage capacity of each table is approximately (2.sup.24) bytes. Accordingly, an extremely large capacity of approximately 50 megabytes, i.e., (2.sup.24) 3, is needed for all of the respective HSL tables, which is impractical."}
{"text": "1. Field of Invention\nThis invention relates to the field of tools and more specifically to a wrench with a socket and boss type connecting means.\n2. Description of Related Art\nWrenches must conform to specific parts and the space available. A typical professional or consumer will have multiple wrenches of varying head configurations. In addition to carrying multiple types of wrenches, professionals and consumers often have multiple wrenches having the same wrench head with a varying body or handle shape. Despite the multiple wrenches available on the market, however, a typical professional or consumer often will not have an optimal wrench available for a specific job. Having the optimal wrench available adds efficiency to a project and relieves physical and mental stress.\nWrenches usually have a forged construction, relying on molds for manufacture. The cost associated with making various molds is one factor contributing to the lack of variation of wrenches available in the market. Wrenches also undergo treatments to ensure they withstand the pressures and forces experienced when performing a given job.\nOne disadvantage to needing multiple wrenches is the large storage space required. A professional may require multiple toolboxes when transporting wrenches to a job site, or may need to choose a limited number of wrenches. Often professionals need to choose a limited set of wrenches without having knowledge of the job site or what specific obstacles they will encounter.\nSelecting wrenches and limiting a working collection, whether for a specific job or in whole, also limits the number of operations that a user can perform with the tools available. Wrenches are specialized tools, and require a different wrench, or wrench head, to perform each task. Workers are limited to a finite number of wrenching jobs based on the wrenches available.\nA further disadvantage to needing multiple wrenches is the weight of transporting even a portion of the wrenches needed. Wrenches are usually solid to withstand the forces generated when tightening or loosening components and are therefore heavy. Carrying multiple wrenches unnecessarily increases the physical demands of an already physically intensive job.\nExisting modular wrench systems known in the art have limited positions for which wrench heads may be varied. Modular system handles are constructed as single segments having a single fixed position.\nIt is desirable to expand and exponentially increase the operations that may be performed using one or more of the wrenches known in art.\nIt is further desirable to develop a wrench system that may be manufactured using alternative processes, such as polymer-based 3-D printing, for low-torque or low-force applications, or sintered metal 3-D printing for higher torque applications."}
{"text": "The work routines of many tradesmen today require that tools and materials be brought to the job and moved around the job location; and that some work be performed on a workbench. For work such as heavy maintenance, plumbing or carpentry the tools are heavy and numerous. Supplies such as plaster buckets, paint cans, pipe fittings and the like add weight and heft to the worker's on-the-job materials. Many tradesmen transport their tools and supplies in a pick-up truck, thus having to load and unload these items daily, To meet the tool and supply handling needs of diverse workers and tradesmen, a wide variety of work site utility carriers are used.\nWhen tradesmen have to move tools and material around at the work site, the usual transport means is a dolly typically having four wheels and some kind of handle for pulling and steering. The numerous tools themselves must be organized and quickly available for use, and virtually all tradesmen therefore use a portable, lockable metal toolbox. On occasion it is also necessary to move articles besides tools and supplies around the work site that cannot practically be loaded and carried on a dolly, such as a desk, plywood sheets, or a washing machine. For this task a sturdy hand truck is needed.\nWork sites typically require a work bench or table as well. This item too must therefore be available on the site to the workman wherever the need for a work table arises. All of these worksite utility items—the dolly, the toolbox, the hand truck and the work table—must be not just transportable to the site but also must be readily loadable across the tailgate of the tradesman's pick-up truck and rolled into position in the truck bed.\nSome of the prior art utility carriers usefully combine the functions of a hand truck and a tool cart. Examples are U.S. Pat. No. 2,964,328 issued 13 Dec. 1960 to Muir; and U.S. Pat. No. 5,224,531 issued 6 Jul. 1993 to Blohm which also includes a work surface. Hand trucks with permanently mounted tool carts or boxes are exemplified by U.S. Pat. No. 6,170,839 issued 9 Jan. 2001 to Kizewski which also has a dolly mode and a flat work surface; and by U.S. Pat. No. 5,588,659 issued 31 Dec. 1996 to Boes et al. U.S. Pat. No. 5,419,569 issued 30 May 1995 to Walls shows a hand truck outfitted with a storage container and capability to carry heavy cylinders using straps.\nMany hand trucks include a pivoted or extendable flat surface to provide a work bench. Examples are U.S. Pat. No. 2,784,004 issued 3 Mar. 1957 to Hamrick and U.S. Pat. No. 4,934,718 issued 19 Jun. 1990 to Voegele. The hand truck of U.S. Pat. No. 6,341,788 issued Jan. 29, 2002 to Ciccone also shows a dolly mode that optionally provides an extendable work table. U.S. Pat. No. 6,758,482 issued 6 Jul. 2004 to Stallbaumer slows a hand truck-dolly combination with detachable handles, work surface and tool tub that may be placed in a raised or lowered mode.\nSome hand trucks are convertible by various means to a fuller work bench. U.S. Pat. No. 5,642,895 issued Jul. 1, 1997 to Wunder shows a hand truck with retractable legs enabling conversion into a work table. U.S. Pat. No. 5,203,815 issued 20 Apr. 1995 to Miller is a conventional hand truck-dolly that mounts a detachable work table top. U.S. Pat. No. 5,957,472 issued 28 Sep. 1999 to Borgatti is a collapsible hand truck and machine stand. U.S. Pat. No. 6,530,583 issued 11 Mar. 2003 to Mueller shows a dolly for wheeling a power saw attached to a support that swings to the horizontal to provide a table saw when the dolly is tipped upright.\nThe folding work bench/tool table of US patent application publication 2002/0179181 of Dec. 5, 2002 by Murphy combines a dolly frame that mounts a removable tool box, A work table is hinged to two rear legs pivotally mounted on the dolly rear castor frames and two front legs hinged at the dolly frame midsection. The table is held by lockable braces either in an upright mode or a lowered mode in which the table secures the tools.\nWhile the preceding prior art work site utility carriers each offer some versatilities and multiple uses, none provide the full functionalities of a hand truck, a retractable work bench, a safe place to store tools, and a dolly. Moreover, none provide a mechanism specifically adapted to allow the utility carrier to be loaded readily by a tradesman into and out of the bed of a pick-up truck."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for transferring path or routing information of internetworking equipment, such as a router, brouter (which serves both as a bridge and a router), etc., which is connected to communication circuits and communicates path or routing information between networks. Throughout the following description and claims, the terms \"path\" and \"routing\" are used interchangeably.\n2. Description of the Related Art\nConventionally, internetworking equipment (e.g., routers) of this type are connected to communication circuits, such as public networks, ISDNs (integrated services digital networks), etc. In order to enable data communication between the networks, the routers carry out propagation and control of path information (routing information) in the networks. Protocols for the propagation of the routing information include dynamic routing protocols, such as an RIP (routing information protocol). The routers request periodical transmission of all pieces of routing information according to the RIP, to thereby exchange the routing information between the networks.\nThis routing method requires a function (hereinafter referred to as \"periodical update function\") to make a simple path calculation from information for routers to execute a gateway operation for data next and from metric information indicative of the number of the routers (including a local one) and transmit all pieces of routing information periodically (e.g., once in 30 seconds), and a function (hereinafter referred to as \"age-out function\") to erase the router information and metric information when there is no data input for a predetermined period of time. Besides the periodical update function, the routing method covers a function (hereinafter referred to as \"triggered update function\") to transmit only changed routing information, if any, immediately.\nWith use of the RIP, according to the routing information method described above, all pieces of routing information need not be set statically, so that change of the routing information attributable to disconnection of the path can be automatically tackled with speed. In some cases, the routers based on this routing method may carry out routing by utilizing a timer-charged circuit, such as an ISDN. Preferably, in these cases, the connection time for the circuit should be minimized so that the circuit can be internetworked only when necessary.\nWhen periodical updating is carried out, the circuit is frequently internetworked for periodical transmission, and sometimes may be left internetworked. In the case where the routing is effected with use of the timer-charged ISDN, therefore, a method is conventionally used such that the routing information is statically set to ensure practical costs.\nIn general, the routing information can be set statically in a very simple network arrangement (including one or two networks). In the least intricate network arrangement (including three or more networks), however, the routing information cannot be easily set without good knowledge of networks.\nConventionally, therefore, routing protocols are not expected to use the periodical update function. In this case, the purpose of application requires the circuit to be disconnected when there is no communication, and information is held without time-out control. Further, the routing protocols are bound to be worked only in specific circuits between specific routers.\nAccording to this routing method, however, the information can be transmitted only a specific number of times, so that it is expected to be delivered correctly to destination routers without fail during the transmission. Once the information is subjected to mismatching, moreover, it cannot be restored with ease, so that the utmost care should be taken to avoid such mismatching. In case of mismatching, furthermore, the information must be able to be restored.\nWhen the routing information is changed, triggered updating is normally effected such that the information is transmitted immediately. In the case of an interface which has no broadcast function (broadcast-type general message receiving function), e.g., an interface such that the partners of internetworking are changed by the source of transmission using an ISDN, however, a triggered update is transmitted to all the partners, so that the routers are subjected to a substantial load. In the case where the number of partners of transmission is large, in particular, the load on the routers is a great problem."}
{"text": "1. Field of the Invention\nThe invention relates to a motion monitor based on detection of modal noise produced by minute motions in a single- or multimode-optical fiber illuminated by a coherent or partially coherent light source, and more particularly to such a motion monitor for detecting body motions due to breaths, heartbeats, and determining breathing rate and heart rate.\n2. Discussion of Background\nA significant cause of death in infants (birth to about 2 years) is \"Crib Death\" or Sudden Infant Death Syndrome (SIDS). Medical authorities generally agree that some infants simply stop breathing during sleep (apnea) or their heart rates fall dangerously low (bradycardia) and that death from these conditions can be prevented if the condition is detected and help is provided within a short time (one or two minutes) by trained personnel or parents.\nApnea monitors already exist, but their cost creates an affordability problem that limits their use. Also, the existing monitors tend to be unreliable, typically having high false positive rates. Also, existing monitors are obtrusive and difficult to use because the subject must wear a belt or other device connected to the monitor. This further limits their use.\nOf interest in the patent literature are U.S. Pat. No. 4,297,684 to Butter, U.S. Pat. No. 4,843,233 to Jeunhomme and U.S. Pat. No. 4,854,706 to Claus et al.\nButter discloses a fiber optic intruder alarm system for protecting the perimeter of an area utilizing a multimode optic fiber as a deformable sensing element, wherein a length of multimode optic fiber is buried in the ground, coherent light is directed through the buried optic fiber and the speckled output light pattern of the fiber is rectified, integrated and threshold detected to determine the presence of an intruder.\nJeunhomme discloses a fiber optic vibration detecting device, but with no means for quantifying the output of the optic fiber, and suggests that observed vibrations can be used to trigger an alarm and/or lighting in the case of a breathing deficiency of a monitored patient.\nClaus et al disclose a modal domain optical fiber sensor for vibration monitoring and as a mechanical motion sensor. Claus et al teach filtering the speckle pattern at the output of an optical fiber to a single lobe, termed a single speckle or a single correlation cell elsewhere in the scientific literature, and finds application to monitoring force on mechanical structures, such as a cantilever beam.\nHowever, in the prior art, how to reliably derive information relating to breaths, heartbeats, breathing rate, or heart rate of a monitored person is not well documented, and a need exists for such a monitor which is both inexpensive and reliable."}
{"text": "1. Technical Field\nThe present invention relates to a method and system for data processing in general, and in particular to a method and system for providing automatic notification. Still more particularly, the present invention relates to a method and system for providing automatic notification of an end of lease of a computer and its location.\n2. Description of the Prior Art\nMore and more companies prefer computer leasing over purchasing. The lease term of a typical computer lease usually lasts approximately two, three, or four years. At the end of a computer lease, a company needs to return all the lease computers; otherwise, the company must continue to pay at the lease rate for any lease computers that have not been returned. The task of finding all the lease computers is usually quite straightforward when there are only a few lease computers. However, when there is a large number of lease computers, such as several thousand units, installed at various locations within a company, then the task of finding all the lease computers becomes daunting.\nA database can be used to keep track of information such as whether or not a computer lease for a specific computer is about to expire, along with the location of the computer according to its first installation. However, computers tend to be relocated many times within a huge organization after their first installation, and yet the database is usually not updated after the relocation of computers. As a result, the database tracking method becomes very impractical. Consequently, it would be desirable to provide an improved method for providing automatic notification when a computer lease is about to expire, and for providing the location of the computer to which the computer lease belongs."}
{"text": "In recent years, along with the rapid proliferation of the Internet, there has been a demand for an increase in capacity of a communication traffic. In response to this demand, an increase in transmission speed of a single channel per system and an expansion of the number of channels in accordance with the wavelength division multiplexing (WDM) are progressed. In the WDM, plural optical signals assigned to different carrier wavelengths (channels) can be transmitted simultaneously and a communication capacity can be increased according to the number of channels. Further, in the WDM, the respective channel wavelengths are sufficiently separated from one another. For example, the communication capacity reaches 1 terabit/sec by transmitting 100 channels using one common optical fiber with the modulation of 10 gigabit/sec per channel.\nIn the middle and long-distance optical communications in these years, a C band (1530 to 1570 nm) that can be amplified by an optical fiber amplifier (EDFA: erbium-doped fiber amplifier) is used widely. In the WDM, for example, with respect to 100 channels used in optical communications, 100 types of laser devices corresponding to the respective channels are required. Therefore, costs for inventory control, inventory count, and the like of the laser devices are increased. In order to solve this problem, in the middle and long-distance communications, there is a demand for a wavelength tunable laser device that can cover the entire C band singly.\nAn example of the light source of the wavelength tunable laser device includes a light source including a semiconductor optical amplifier (SOA), a wavelength tunable filter, and a reflector (external resonator type wavelength tunable laser or the like). The wavelength tunable filter controls the oscillation wavelength of the laser by changing the refractive index of a filter portion (the refractive index of a material of a filter portion). Examples of the publicly known document in which a wavelength tunable laser device provided with such a light source is described include Patent Documents 1 to 4."}
{"text": "The machines used for these works have a particular morphology. They are generally equipped with stocks or a torsion bar system making it possible to position the profile section to be machined under a broaching tool and to drive it with the movements necessary to the milling operations that make it possible to produce the part.\nThese movements are combinations of longitudinal displacements and rotations about the long axis of the part.\nTo machine all the faces of a profile section in order to obtain a finished part, it is necessary to work in at least 2 phases having different lays. The first phase is for machining a part of the surfaces by taking reference on two rough faces. The second is for machining the remaining rough faces, after reversing the part and calibrating on the previously machined faces.\nThis part reversal and recalibration operation is performed manually by the operator of the machine. The invention makes it possible to perform this operation automatically. It thus eliminates the fatigue associated with this operation for the operator and the risk of poor positioning of the part after reversal. It thus allows for a significant productivity gain."}
{"text": "Capacitive touch control technologies may sense whether there is a user's finger or other conductive object by a detection capacitor, and a corresponding function can be implemented simply by a light touch with a finger. Currently, touch screens have been widely applied to electronic devices, to implement touch control operations.\nIn the process of conceiving the present disclosure, the inventor finds that there are at least the following problems in the existing technologies:\n(1) Existing electronic devices support only one-dimensional input, but cannot support multi-dimensional input, and pressure detection is not sufficiently accurate.\n(2) A simple finger touch control function is only applicable to a limited number of application scenarios. This does not adapt to the development trend of technologies of integrating multiple application functions in one electronic device."}
{"text": "This invention relates to a method of depositing a high-emissivity carbon coating on a substrate by reactive RF sputtering. More specifically, the invention relates to a method of depositing a carbide-carbon coating on a substrate by reactive RF sputtering having a composition gradient such that the coating is composed primarily of a carbide compound near the substrate and primarily of carbon near the exterior surface.\nThermoelectric generators are frequently used as an energy source in aerospace applications. Temperature regulation in these generators is a critical problem for both orbital satellites and deep space probes. A high-emissivity coating on the metallic thermoelectric generator casings would improve the performance of satellite thermoelectric systems by radiating large quantities of heat while operating at low temperatures. The high-emissivity coating must be capable of withstanding the stresses imposed during launch and the elevated temperatures and vacuum conditions encountered during operation.\nCarbon has many of these properties and would serve well as a high-emissivity coating. However, carbon does not adhere well to the metals often used in thermoelectric generator casings, namely iridium and aluminum. Attempts have been made to blend carbon with an epoxy and paint the mixture onto the casings, but this process adds undesirable weight to the generator. A more promising method is to sputter deposit a coating with a composition gradient such that the surface of the coating near the substrate is rich in material that adheres well to the substrate while the exterior surface of the coating is rich in carbon. Previous methods of deposition have involved rotating a substrate between two targets of different materials and gradually changing the deposition parameters on the two targets to obtain the desired composition gradient. This method is too complicated to be practical for the massive parts of the generator which must be coated."}
{"text": "In the field of image processing, neighborhood masks are commonly employed when it is necessary to consider the environment of a pixel to be processed. A neighborhood mask is defined as a selection area in the image. With the changes in image processing, the neighborhood masks are becoming increasingly extensive. At the same time, the research and development and manufacturing costs of architectures supporting image processing place major constraints on their development in terms of silicon surface area and electricity consumption, notably when these architectures are designed for embedded applications intended for the general public. The issue in this field is therefore to have open-ended and flexible computation architectures which, given the same circuit, are capable of supporting the existing processing operations as well as those to come whose computational complexity will be greater. One of the known solutions for meeting these needs consists in using modular architectures, based on stream processing units each comprising a set of neighborhood processors operating in parallel on neighboring data. A first exemplary architecture is described in the patent application EP 08 05369. Each processing unit comprises a set of processors and a storage unit containing all the neighborhoods accessible to the processors. However, a processor can access only the neighborhood which is assigned to it. Furthermore, the dimension of the storage unit is fixed. Consequently, the size of the neighborhoods accessible to the processors is limited to that provided at the time of the design of the circuit. A second exemplary architecture is proposed by the Korea Advanced Institute of Science and Technology and described in K. Kim et al., “A 125 GOPS 583 mW Network-on-Chip based parallel processor with bio-inspired visual attention engine”, IEEE Journal of Solid-State Circuits, vol. 44, no. 1, January 2009. Each processing unit comprises a set of processors and a local memory. Each processor has full access to the local memory of the processing unit to which it belongs. The local memories of different processing units can communicate with one another and each processor can communicate with left and right neighboring processors, including if a neighboring processor is located in another processing unit. These communications between local memories and between processors make it possible to extend the size of the neighborhoods accessible to the processors. The architecture is, however, ill suited to such a use: on the one hand, the communication from processor to left/right neighboring processor between processing units allows only one datum to be exchanged per cycle, which in practice limits the extension of the size of the neighborhoods. Also, the communication from local memory to local memory entails overwriting at least a portion of the data present in one of the local memories, which constrains the order in which the data must be processed."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a process for the modeling and control of production integrated-processing equipment (PIPE). It further relates to a class of PIPEs which have operations sequences. It also relates to the use of PIPE-specific models to control material movement within the PIPE and loadings of the PIPE. It also relates to the use of such PIPE-specific models to analyze and extrapolate PIPE performance. It also relates to the use of PIPE-specific models to evaluate and develop PIPE designs and PIPE control schemes. The invention further relates to a computer integrated manufacturing (CIM) system that includes a PIPE-specific model for automatic PIPE control and to the use of PIPE-specific models in the design of CIM systems.\n2. Description of the Prior Art:\nMany manufacturing plants or factories are using or introducing highly-complex production integrated-processing equipment (PIPE). This integrated equipment is also known as modular equipment or cluster tools. PIPEs are capable of performing multiple, independent processing steps without leaving a controlled environment. By performing sequences of processing steps in one machine, substantial yield improvements over conventional processing equipment may result. This automated processing reduces contamination and unnecessary human handling. The complexity of this integrated equipment, however, introduces severe problems in analysis of the PIPE performance.\nPrior to this invention, PIPE manufacturers and users have taken one of two approaches to understand PIPE performance: spreadsheets or physical experimentation. Spreadsheets provide a basis for approximate analytic calculations. While such an approach results in computer calculations that provide impressive-looking numeric results, little or no effort is expended in establishing the underlying accuracy of such spreadsheet calculations. A major area of approximation is the description of internal movement through the PIPE. This internal material movement is complex, and poorly understood. One approximation in spreadsheet calculations has been to surrender to this complexity and describe the operations sequence or path through the PIPE as being random.\nThe second approach taken by PIPE manufacturers and users has been physical experimentation. There is clearly a need for performance analysis, especially with regard to throughput and cycle time, since PIPEs, in general, are multi-million dollar pieces of equipment. A single physical experiment can cost tens-of-thousands of dollars and provide only one data point. The need for performance analysis has compelled manufacturers and users to perform such experiments in order to improve on the accuracy of spreadsheets. Millions of dollars may be spent doing performance analysis by this method. Any answers obtained by this method are case-specific, and any change in physical design, operating conditions, or recipes will negate the results and require repeating the effort.\nPerformance Analysis and PIPE Productivity\nOne of the goals of performance analysis is to calculate the cycle time of a standard manufacturing unit, i.e., part, lot, or cassette, as a function of process recipe. The process recipe is the set of processing steps to be performed by the PIPE. Despite the processing advantages provided by PIPEs, realistic calculations of throughput and capital cost per part impact the user's choice of a PIPE for a particular process step or steps. The calculation of throughput of a PIPE workstation is based primarily on the PIPE's cycle time.\nTable 1 demonstrates the impact of cycle time on the capital cost per part. In this example, 25-wafer cassettes are being processed by a PIPE. The PIPE in this example has a modest capital cost of $1 mm which is amortized over 5 years. The capital cost per part is a direct function of the cycle time of the cassette in the PIPE. A variation of a factor of four in cycle time results in a variation of a factor of four in the capital cost per part. Clearly the accurate calculation of PIPE cycle time is crucial information in making decisions on equipment selection and numbers of machines needed at a workstation. Further, the total operating cost per wafer may be twice the capital cost.\nTABLE 1 ______________________________________ PIPE PRODUCTIVITY vs CYCLE TIME ASSUMPTIONS: PIPE CAPITAL COST = $1,000,000 5 YEAR AMORTIZATION ($200,000/yr) 325 WORK DAYS ($615 in amortization/d) 20 PROCESSING HOURS PER DAY CASE 1 CASE 2 ______________________________________ CYCLE TIME 0.5 HRS 2 HRS # CASSETTES/d 40 10 COST/CASSETTES $15.38 $61.50 COST/Part $0.62 $2.45 ______________________________________"}
{"text": "1. Field of the Invention\nThe present invention relates to methods and apparatus in the field of cardiotomy perfusion, and in particular, it is directed to cardiotomy reservoirs which are used in extracorporeal blood circuits.\n2. The Prior Art\nThere are a wide variety of medical and surgical procedures in which a patient's blood is conducted extracorporeally through a blood circuit for treatment and thereafter returned to the patient. Typical extracorporeal blood circuits include, for example, cardiopulmonary bypass and extracorporeal hemodialysis.\nEffective circulation and control of a patient's blood are essential to successful heart surgery. In modern surgical procedures, blood is circulatd through the patient mechanically by establishing a cardiopulmonary bypass circuit in which the blood is forced through a bloodline or tube with a blood pump. This circuit removes blood from the patient, oxygenates the blood, and then returns the blood to the patient.\nIn addition, blood which accumulates in the patient's pericardium or heat cavities during the course of surgery is transported away from the surgical site. This is generally accomplished by means of a \"cardiotomy sucker\" which aspirates the blood from the areas of accumulation. During an operation of only thirty minutes, as much as 250,000 cubic centimeters of blood may be aspirated from the patient through the \"cardiotomy sucker.\" Consequently, unless it can be successfully returned to the patient, this aspirated blood must be replaced by blood from other sources. In attempting to return this aspirated blood to the patient, however, one encounters of number of obstacles.\nBlood removed from the surgical site through a \"cardiotomy sucker\" invariably contains tissue fragments, blood clots and other debris which prevent the blood from being returned directly to the patient through the cardiopulmonary bypass circuit. This difficulty can generally be minized by allowing the blood to pass through filters which remove all particles larger than a predetermined size, for example, twenty (20) microns. However, this filter will sometimes become occluded and must, consequently, be replaced.\nA somewhat more difficult obstacle to recycling aspirated blood is due to the fact that blood which is removed from the surgical site through a \"cardiotomy sucker\" has a substantial amount of air bubbles suspended within it, thereby resulting in a foam-like solution. Oxygenators which are used in the cardiopulmonary bypass circuit have traditionally had some limited defoaming capability, because small air bubbles remain in the blood after the oxygenating process. However, aspirated blood differs greatly from oxygenator blood for which the oxygenator defoamer was designed. Although bubbles are present in both oxygenator and aspirated blood, the bubbles contained in oxygenator blood are typically within the same, specific micron range as the bubbles produced by the oxygenator. Aspirated blood, on the other hand, contains bubbles with no such size uniformity, and the quantity of bubbles may vary from a steady stream to a froth. As a result, defoaming aspirated blood is significantly more difficult than defoaming oxygenator blood and has provided a significant challenge to those skilled in the art.\nIn order to more effectively prepare aspirated blood for return to the patient, a device known as a cardiotomy reservoir has been used. Cardiotomy reservoirs currently in use have the capability of holding large volumes of blood separate from the main extracorporeal cardiopulmonary bypass circuit. Many of the cardiotomy reservoirs currently in use also include internal filters which separate out the above-mentioned debris. Typical prior art cardiotomy reservoirs also contain a small defoaming element which defoams the blood as it passes through the cardiotomy reservoir. After the blood has been so treated, it is then recycled to the patient through the cardiopulmonary bypass circuit.\nAlthough the cardiotomy reservoirs presently in use have greatly improved the quality of the blood which is returned to the patient, current cardiotomy reservoirs have proven inadequate in many respects. Current cardiotomy reservoirs use a defoaming substance which is composed of the same material used in the defoamer within the oxygenator. This defoamer, which is a uniform single density defoamer, cannot effectively defoam aspirated blood where there are a wide variety of sizes and quantities of air bubbles suspended within the blood.\nAdditionally, the design of present cardiotomy reservoirs allows blood to pass quite rapidly through the reservoir. That blood is, therefore, not in contact with the defoaming element for any significant length of time and this results in inadequate defoaming of the blood. Even if the outflow of blood from current cardiotomy reservoirs is intentionally interrupted, significant volumes of blood are held in the lower portion of the cardiotomy reservoir and are not in contact with the defoamer element at all. Consequently, the temporary interruption of blood outflow from present cardiotomy reservoirs does significantly improve their defoaming capability.\nFinally, cardiotomy reservoirs which include internal filters occasionally clog due to excessive debris and/or blood clots resulting from inadequate heparinization, accelerated metabolism of the heparin, or blood coagulopathy. When the internal filtr of the cardiotomy reservoir becomes clogged, the entire cardiotomy reservoir must be changed. This interrupts and disrupts the cardiotomy reservoir circuit. Equally important, the blood trapped within the clogged cardiotomy reservoir cannot always be recovered and returned to the extracorporeal circuit; hence, significant volumes of blood must sometimes be discarded. Even if the trapped blood is recovered, it will often not be used because there is an insufficient guarantee of sterility. In addition, there is often even no attempt to recover this trapped blood, because the time needed to recover the blood to the extracorporeal circuit may be too lengthy. Each of the above-mentioned situations adds significantly and dangerously to the normal risk involved in heart surgery.\nAccordingly, it would be a significant improvement in the art to provide a cardiotomy reservoir having a defoamer which is specifically designed to handle blood containing suspended air bubbles of many sizes, such as blood aspirated during heart surgery, Additionally, it would be an improvement in the art to provide a cardiotomy reservoir in which the defoamer is always in substantial contact with the blood which is located within the reservoir. Further, it would be an improvement in the art to provide a cardiotomy reservoir which would allow for the rapid replacement of a clogged filter, without compromising asepsis or blood recoverability. Finally, it would be an improvement in the art to provide a method for defoaming aspirated blood which adequately separates the air from the blood while minimizing the risk to the patient incident to the defoaming process. Such a novel cardiotomy reservoir apparatus and method is described and claimed herein."}
{"text": "1. Field of the Invention\nThe present invention relates to a folding ladder tree stand and more particularly pertains to providing a hunter an elevated stationary support for hunting with a folding ladder tree stand.\n2. Description of the Prior Art\nThe use of tree stands is known in the prior art. More specifically, tree stands heretofore devised and utilized for the purpose of providing a stationary support for hunting are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nBy way of example, U.S. Pat. No. 4,129,198 to Hunter discloses a pole seat and ladder. U.S. Pat. No. 4,742,888 to Amacker discloses a folding ladder stand. U.S. Pat. No. 5,242,030 to Lobozzo discloses a deer stand. U.S. Pat. No. 5,253,732 to Daniels discloses a portable folding tree stand. U.S. Pat. No. 5,267,632 to Mintz discloses a folding deer stand apparatus.\nWhile these devices fulfill their respective, particular objective and requirements, the aforementioned patents do not describe a folding ladder tree stand that provides an elevated stationary support for hunting, provides a way to readily access this elevated stationary support, and provides a structure that can be readily placed in an operable configuration for use or a stowed configuration for transport.\nIn this respect, the folding ladder tree stand according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in doing so provides an apparatus primarily developed for the purpose of providing a hunter an elevated stationary support for hunting.\nTherefore, it can be appreciated that there exists a continuing need for new and improved folding ladder tree stand which can be used for providing a hunter an elevated stationary support for hunting. In this regard, the present invention substantially fulfills this need."}
{"text": "1. Field of the Invention\nThe present invention relates to a quantitative dispenser, and in particular, to a quantitative dispenser which comprises a divider a with scale or scale-like means that is adapted to dispense a constant amount of powdery material such as powdered milk stored in a volumetric body.\n2. Description of the Prior Art\nThere are many powdery or grain-type materials in our daily lives, for example, powder detergent, nourishing drug pills, healthy medicines, powdered milk, and the like, each of which has to be administrated or used at its own specified dosage. Each powder thereby needs its own measuring means, i.e., cup, spoon or the like for dispensing the powder from its container. These measuring means thus have various sizes depending on the kind of products and various serving conditions, and hence most of them cannot be used in common. This results in production of a lot of garbage at the end of their use, which is a waste of natural resources. In addition, big cans containing a large volume of powdery material are inconvenient to carry. This often causes trouble for their users, such as the use of many small cans or removing fractions with paper or the like.\nAs an example, currently, there is a powdered milk box which consists of 3 or 4 cylindrical compartments for containing powdered milk. These compartments are engaged with each other through screw threads. At the top of the box, there is a funnel with a screw cap on its central outlet. That funnel can be engaged with each compartment through, again, the screw thread. When dispensing, a screw cap on the outlet of the funnel and turning the box upside down, allows powdered milk in the first compartment to be poured into a milk bottle or the like.\nHowever, there are several disadvantages for this kind of dispensing box: for example, a certain amount of powdered milk has to be loaded in each compartment before dispensing. The amount of powdered milk can not be varied according to different situations so that the amount of milk dispensed may be too much or too little. Also, the number of usages for the container is limited to the number of compartments.\nAccordingly, there exists a need to provide a more economically efficient and more convenient quantitative dispenser in order to overcome the disadvantages of the prior art.\nAccordingly, the objects and advantages of the invention can be achieved by providing a quantitative dispenser which will be described hereinafter in detail with reference to its preferred embodiment."}
{"text": "1. Field of the Invention\nThe present invention relates to electric lamps and more particularly, but not by way of limitation, to a lamp having a treated lense to reduce peripheral glare primarily for use as automobile or vehicle driving lamps, spotlights or the like.\n2. History of the Prior Art\nThere is a well developed history of utilization of lamps for spotlights or driving lamps where the lamp includes an internal parabolic reflecting surface for collecting the light emission from a bulb or resistive element located near the focus of the parabolic surface, collimating the light into a beam which is directed through the lense or suitable transparent media to illuminate a specific viewing area.\nSince the lense is usually refractive and the parabolic surface is often slightly open to provide a desired beam pattern, there exists peripheral illumination outside of the beam width in which the light is less intense and which is even desirable for many ordinary uses.\nHowever, when such lamps are used as driving lights for vehicles in fog, smoke or rainy conditions, reflection from the peripheral illumination is distractive and actually impairs visibility in the desired viewing area.\nA particularly successful attempt was made to overcome the above disadvantages by the patent to Hulbert, U.S. Pat. No. 3,754,135, issued Aug. 21, 1973, for \"Light Treating Means.\" The Hulbert device teaches the treating of the light by coating specific areas of the lense with a blue translucent pigment whereby the teaching indicates that a \"light mixing\" takes place which serves to change the color spectrum and to reduce glare. While glare is in fact reduced by the Hulbert device, a great amount of peripheral illumination is still present which can and does present some problems when the lamp is used in conditions of fog, smog, rain or on wet surface.\nOn the other hand, by utilizing the blue mixed light, there is an enhancement in color differentiation and object identification within the illuminated beam pattern of light.\nAnother problem exists with the Hulbert device when used as a vehicle head lamp that is that the primary beam is surrounded by a reflected blue halo during use in highly reflective environments which can be distractive to approaching motorists as well as to the driver of the vehicle. Also, since the Hulbert lamp appears blue during daylight conditions as well as when used in a highly reflective environment, its use is prohibited under many state laws in which blue has been adopted for use by law enforcement vehicles only."}
{"text": "1. Field of the Invention\nThis invention generally relates to the measurement of a multiple phase flow. Here, microwave radiation is used to measure the two-phase fluid content of oil field flowlines and also the two-phase fluid content of fluid-saturated porous media.\n2. Background Information\nA multiple fluid mixture has at least one physically distinct and mechanically separable fluid within at least one other physically distinct and mechanically separable fluid in a physical/chemical system. In the case at hand a multiple fluid mixture would be something such as a combination of water and oil.\nSome of the factors that complicate the measurement of the multiple fluid mixture are: whether the mixture is in a static or dynamic state; and whether the method is intrusive or non-intrusive. One simple method of analyzing a multiple fluid mixture involves capturing of a sample volume of the mixture and then analyzing it (or simply allowing the mixture to separate into its component phases). This method, however, is an intrusive one and involves too much time and effort to capture a sample volume. Furthermore, when captured, there is no guarantee that the sample volume is an adequate representation of the entire sample. If so, the calculation for the total concentration of the two components in the mixture will be erroneous. For that reason a non-intrusive method would be preferable.\nMultiple fluid mixtures may also be measured by subjecting a sample (either in a static or dynamic state) to photonic laser, or infra-red radiation. In these cases, a multiple fluid mixture is subjected to the specific radiation and the attenuation of the radiation is used to calculate the component phases. See U.S. Pat. Nos. 4,157,470 to Kotuka et al.; 4,470,647 to Chraphyvy et al.; and 4,394,575 to Nelson; see also V. E. Shrock, \"Radiation Attenuation Techniques in Two-Phase Flow Measurements\", Nat'l ASME/AICHE, Heat Transfer Conference, Minneapolis (1969); and Davis, \"VHF Electrical Measurement of Saturations in Laboratory Corefloods\", SPE 8847 (1980).\nThe other factor that was mentioned above, i.e., whether the mixture is in a static or a dynamic state, would reflect mainly on the particulars of the measurement system. For example, in most applications it is most desirable to measure the mixture when it is in a fluid or flowing state. A common example would be an oil field flowline pipe that may transport liquid. However, there are examples where it may be desirable to measure multiple fluids in a static system. For example, a laboratory situation may be set up in which a sample petroleum formation, containing a porous medium (such as Berea Sandstone), has a two-phase mixture. Both cases could be applications of the same measurement technique.\nThe techniques described above are useful in understanding how fluid concentrations may be measured in a multiple fluid mixture. However, as applied to petroleum applications, i.e., multiple fluid mixtures of water and oil in a flowline or a test core, microwaves work to distinguish between multiple fluids. Some work has already been done in which microwaves are used to determine the concentration of multiple fluids, as outlined below. The devices disclosed in U.S. Pat. Nos. 4,167,736 to Tomlinson; 4,196,385 to Vesterpaard et al.; 4,423,623 to Ho et al.; and 4,503,384 to Nagy et al. all are intrusive into the fluid that they are measuring and are undesirable because of that fact. U.S. Pat. Nos. 4,289,020 and 4,301,400, to Paap, show devices that manipulate the fluid channel through which crude oil flows and are undesirable because they contain \"dead spaces\" where the flow may eddy and throw off the measurement system. U.S. Pat. No. 3,586,971 to Besido shows a device that monitors the moisture content in the sheets of material. It uses a method that detects a resonant frequency shift only. U.S. Pat. No. 3,818,333, to Walker, shows a microwave window and antenna, but does not disclose how measurements are made to determine fluid content. The following patents are issued to Lorne A. Davis et al. U.S. Pat. No. 4,519,982 discloses a device that scans a two-dimensional sample. U.S. Pat. Nos. 4,482,634 and 4,490,676 disclose a method ('634) and an apparatus ('676) for monitoring fluid in a core of material. However, both patents only measure power level, not phase shift. U.S. Pat. Nos. 4,486,714 and 4,543,821 use microwave circuitry that is similar to 4,482,634, but show different methods of fluid delivery to a test core. Two articles relate to some of these patents and are: Haskin and Davis, \"A Comparison of Laboratory Linear and Pattern Flow Chemical Floods Using a Volumetric Linear Scaling Concept for Oil Saturation and Distributions\", SPE 10197 (1981) (see U.S. Pat. No. 4,519,982); Brost and Davis, \"Determination of Oil Saturation and Distributions in Field Cores by a Microwave Spectroscopy\", SPE 10110 (1981) (see U.S. Pat. Nos. 4,482,634 and 4,490,676).\nOther articles disclose the use of microwaves to measure fluid contents. R. W. Parsons was one of the first investigators to disclose a method to determine in situ saturations of multiple fluid phases in a laboratory coreflood using microwave attenuation, see Parsons, \"Microwave Attenuation--A New Tool for Monitoring Situations in Laboratory Flooding Experiments\", Soc. Pet. Eng. Jour. 302 (August, 1975). However, the signal penetration depth in this method is limited to one inch and the frequency is 21.26 Gigahertz. P. E. Baker discusses a similar method (at the end of the Parsons' paper, see page 309). Bentsen and Saeedi show a method that is similar to Parsons, see Bentsen and Saeedi, \"Liquid-Immiscible Displacement in Unconsolidated Porous Media\", Jour. Can. Pet. Tech. 93 (January-March, 1981). Their method measures only microwave attenuation, has a penetration depth of only 1.27 centimeters, and operates on a microwave frequency of 27 Gigahertz. Wasan et al. also show a method similar to Parsons, see Wasan et al., \"Microwave Spectroscopy Analysis of Surfactant/Polymer Flooding: Interrelationships Between Chemical Slug Properties, Coalescence Phenomena, and Tertiary Oil Recovery\", SPE 8327 (1979). For example, Wasan et al. measure microwave attenuation using a frequency of 21.83 Gigahertz. Another similar device is described in Gladfelter and Gupta, \"The Effect of Fractional Flow Hysteresis on Recovery of Tertiary Oil\", Soc. Pet. Eng. Jour. 509 (December, 1980). Here, measurements of multiple fluids are made with microwave attenuation and depth penetration is limited to 1.91 centimeters. An article that reviews all the current methods to measure multiple fluids and petroleum applications is Saraf, \"Methods of In-situ Saturation Determination During Core Tests Involving Multiphase Fluid Flow\", Report 1981-6 Petroleum Recovery Institute (March, 1981).\nThese references mostly use only one measurement technique, i.e., signal attenuation, to determine fluid concentrations. This has some drawbacks, in that: some physical properties are masked in different flow regimes when you use signal attenuation; you may be limited to determining the concentrations of only two fluid phases (with one measurement technique); and you do not have a confirmation of the fluid concentrations when you use a single method."}
{"text": "The present invention relates to a tool that is used in the construction of frames.\nFor purposes of this description, a frame is a structure having multiple members that are fitted together and united. In constructing frames, attachment of the individual members to one another, by welding or similar means, often requires the use of clamps or vises. Even when using such clamps, it is often difficult to secure adjacent members at precise angles.\nFor example, in manufacturing awnings, a fabric cover is fit over a frame typically formed by a segments of metal tubing. These frames may define a simple rectangular volume or a complex geometric volume. In constructing such a frame, each individual member is ordinarily attached to adjacent members by a weld. Clamps are commonly used to secure adjacent members at the proper angle and orientation for welding. However it is still difficult to secure adjacent members at precise angles using clamps. When more than two members meet at a joint, it becomes even more difficult to secure the members at the proper angle and orientation for welding. A series of welds is often needed to secure a joint uniting more than two members. In short, construction of these awning frames is often difficult and time-consuming.\nIt is thus an object of the present invention to provide a tool that facilitates the securing and welding of frame members at precise angles and orientations.\nIt is a further object of the present invention to provide a tool that allows a single weld to be used to secure three or more members at a joint.\nIt is still a further object of the present invention to provide a tool that is easy to use, even in construction of complex frame structures, such as those used in the manufacture of awnings."}
{"text": "1. Field of the Invention\nThe invention relates to signal transfer points (STPs) used in telecommunication networks and, more particularly, to cabinet designs for STPs.\n2. Description of the Related Art\nAdvanced Intelligent Networks (AIN) are modern telephone networks that separate call control from the actual trunk that carries the voice or data transmission of a call. In general, AIN Networks include a system in which the network queries a database to determine how a call should be processed once a set of call digits have been received from a phone. AIN Network generally includes three basic elements. A signal control point (SCP) is a computer controlled database that is for storing customer\"\"s specific information that is used by the network to determine how to route and/or process calls. A signal switching point (SSP) is a digital telephone switch that is operable to communicate with SCPs and to obtain the customer\"\"s specific instructions for processing the call. A signal transfer point (STP) is a packet switch that shuttles messages between the SSPs and the SCPs. All three, namely the SCP, the SSP, and the STP communicate via what is known as out-of-band signaling. Out-of-band signaling generally refers to the signaling that is for controlling a call. In other words, out-of-band signaling does not include the actual call data or voice information. Typically, the three described systems utilize a signaling system No. 7 (SS7) protocol or a variant thereof.\nWhen a call is dialed, the SSP creates a query to the database within the SCP to find out how a call should be processed. The query is passed via out-of-band signaling through at least one STP to the SCP. The SCP interprets the query based on criteria within its databases and based upon information provided by the SSP. Once the SCP retrieves and transmits a return message through the STPs to the SSPf SSP may properly process the call within the network.\nThe STP is, in essence, a router for the SS7 network. It relays messages through the network but does not originate them. It is similar to a voice switch except that it only routes control signals. STPs also serve as gateway devices that may convert messages that flow between dissimilar systems. For example, a STP within the US may provide conversions between ANSI SS7 and ITU-T SS7. STPs may also provide other functionality including various security functions and message filtering wherein messages of a non-conforming nature are not passed through from one point to another. Other STP functions include traffic monitoring for billing purposes and for developing usage statistics.\nCompact PCI is an adaptation of the Peripheral Component Interconnect (PCI) specification for industrial and/or embedded applications requiring a more robust mechanism form factor than desktop PCI. Compact PCI uses industry standard mechanical components and high performance connector technologies to provide a system optimized for rugged applications. Compact PCI is electrically compatible with the PCI specification allowing low cost PCI chip sets to be used in a mechanical form factor suited for rugged environments. In general, Compact PCI is a rugged variation of the PCI bus that was designed by PC makers wanting to create industridi grade PCs. The motivations for the design were to put more PCI cards into one PC (eight instead of four) and to make the resulting PC more rugged, and better able to withstand shaking and rough environmental conditions.\nThe physical configuration of the hardware conforms to the Eurocard (VMR-style) standard. The cards are identical to VME cards in size but differ in that they use a high-density 2-millimeter pin and socket connector (contact space) for interfacing the card to a passive backplane. Typically, the Compact PCI comes in a rugged 3U or 6U Eurocard form factor and has either a 32 or 64 bit data bus with transfer rates of up to 528 megabits per second.\nAs mentioned before, Compact PCI typically includes a passive backplane. Each shelf on the front side of a Compact PCI system includes slots for eight cards. Traditionally, a controller goes into slot one leaving seven slots for other specialized peripheral cards. Accordingly, cabinets made to house Compact PCI equipment typically provides cooling only to a front portion of the cabinet. The xe2x80x9cfront portionxe2x80x9d is the portion that is in front of the backplane.\nTraditionally, active cards go in the front side while passive circuitry goes on the backside. In a typical configuration, a Compact PCI card having active elements includes pin outs that are passed through the backplane to passive transition modules on the back side. The transition modules terminate the signals into connectors wherein that may readily be accessed from behind. Because the circuitry on the back side of the backplane is passive, there has been no need hereto before to provide cooling or thermal management for the back side of the Compact PCI cabinet in a telecommunications environment.\nA current design goal for the present invention is to utilize Compact PCI in the design of modern telecommunication systems. As systems become more complex, however, Compact PCI systems are being pushed to the limit in terms of capacity. A typical cabinet for a Compact PCI system may have several shelves of cards all generating tremendous amounts of heat energy for a cabinet of that size. Accordingly, thermal management issues may become acute. As telecommunication networks continue to increase in sophistication, there will be an ever-increasing need to add more circuitry within the Compact PCI cabinets. Unfortunately, however, this also means that more space must be found to provide the additional card slots. Greater thermal management must be achieved to maintain temperatures at an acceptable level if new cards are added. What is needed, therefore, is a cabinet architecture that provides for greater circuit capacity and yet accommodates an efficient thermal management system therefor, preferably within the context of a Compact PCI-compliant equipment rack.\nThe Compact PCI cabinet includes an inlet port for receiving external air, a set of blowers for drawing the air into the cabinet and for expelling the air from an exhaust port of the cabinet after having cooled internal active circuits, and a plurality of louvers and openings that create thermal convection within a back portion of the Compact PCI cabinet to cool active components coupled to the back side of the backplane of the Compact PCI system. More specifically, a bottom panel, lining one wall of the chamber that receives the external air through the input port, includes a plurality of louvers that allow some of the air to be conducted into the back portion of the Compact PCI cabinet. Additionally, a plurality of perforated air openings located approximately at the bottom and back end of a side panel of the Compact PCI cabinet also is included to receive and conduct air from the external environment. Finally, the invention includes forming openings between a plurality of card guides located approximately at the top of the back side of the back portion of the Compact PCI cabinet wherein air may be conducted into the back portion of a Compact PCI cabinet or out of the back portion of the Compact PCI cabinet. The air is conducted here, from and into the exhaust from the blowers.\nBecause exhaust from the blowers is being partially directed into the back portion of the Compact PCI cabinet, that exhaust must be cooled before it is used to cool active components in the rear portion of the cabinet. Accordingly, a pin field is placed in the air path between the blowers and the openings through which the air may be conducted into the back portion of the Compact PCI cabinet. The pin field is coupled to a heat sink that is for extracting heat energy from the pins. Accordingly, heat is convected away from the exhaust into the pins and up into the heat sink where it is radiated outwardly from the cabinet.\nFinally, the invention includes placing a thermistor in the exhaust path so that the fan speed of the blowers maybe adjusted according to the temperature of the exhaust after passing through the pin field. Thus, acoustic concerns are addressed by minimizing fan speed and associated noise levels so as to operate the blowers at nearly a minimal level required for maintaining the internal temperatures within the PCI cabinet at an acceptable level.\nIn another embodiment of the invention, the louvered panel at the bottom of the Compact PCI cabinet is installed at an angle relative to vertical so that a greater amount of the air received at the inlet port is conducted into the back portion of the Compact PCI cabinet and so that the louvers create lower amounts of air resistance to that air which is being conducted into the back portion of the Compact PCI cabinet."}
{"text": "Wireless and wired sensors are used in a wide variety of industries. For example, wireless pressure sensors are routinely used in automobiles and other vehicles, such as to monitor the pressure of tires on the vehicles. Also, the tires of a vehicle can often suffer from excessive temperatures, such as due to overloading of the vehicle or high temperatures in some regions of the world. Wireless temperature sensors can be used in a vehicle to monitor the temperatures of its tires.\nTire monitoring systems routinely need to be low in cost and use sensors that are wireless, small in size, and highly reliable with low drift. Existing tire pressure monitoring systems often use piezoresistive pressure sensors. These types of monitoring systems often require the use of radio frequency (“RF”) electronics as well as signal conditioning and processing circuitry. These types of monitoring systems also often require the use of batteries to power the sensors and related electronics and circuitry. However, these types of monitoring systems are often not attractive due to the weigh of the systems and due to the limited lifetime of their batteries, which can be discharged much faster in the harsh environment of vehicle tires."}
{"text": "This invention relates to a separation system and process as illustrated in FIG. 1 useful for the product produced by Armstrong method as disclosed and claimed in U.S. Pat. Nos. 5,779,761; 5,958,106 and 6,409,797, the disclosures of each and every one of the above-captioned patents are incorporated by reference."}
{"text": "Multi-focal type optical lenses having a plurality of focal points are known from the past, as one type of ophthalmic lenses used for a human eye optical system. For example, with contact lenses used as corrective optical elements for refractive error, alternative optical elements after lens extraction or the like with the optical system of the human eye, or with intraocular lenses used for insertion in the human eye, by applying multi-focal lenses, it is possible to compensate for the decrease or loss of accommodation function of eye in the human body.\nParticularly in recent years, there is an increase in people continuing to use contact lenses even when they reach the age of having presbyopia. People with presbyopia have a decrease in focus accommodation function, so a symptom appears of having difficulty focusing on nearby items. Thus, multi-focal contact lenses which can also focus on nearby objects become necessary for presbyopia patients. Also, for patients who have undergone cataract surgery, the lens which is in charge of the adjustment function is removed, so even if an intraocular lens is inserted as a replacement, the symptom of difficulty seeing close up remains. A multi-focal function that offers a plurality of focal points is necessary for that intraocular lens as well. Thus, there is a great increase in the need for multi-focal lenses reflecting the aging society of recent years.\nHowever, as a method for realizing this multi-focal lens, examples are known of a refraction type multi-focal lens for which a plurality of focal points are formed based on the principle of refraction, and of a diffractive type multi-focal lens for which a plurality of focal points are formed based on the principle of diffraction. With the latter diffractive type multi-focal lens (diffractive multi-focal lens), equipped are a plurality of diffractive structures formed in concentric circle formed on the optical part of the lens, and a plurality of focal points are given by the mutual interference effect of light waves that passed through the plurality of diffractive structures (zones). Thus, compared to the refraction type lens with which a focal point is given by the refraction effect of light waves at a refracting surface comprising boundary surfaces with different refractive indexes, with the diffractive type multi-focal lens, there are advantages such as being able to set a high lens power while inhibiting an increase in lens thickness.\nTypically, the diffractive multi-focal lens has a diffractive structure by which the diffractive zone pitch gradually becomes smaller as it goes from the lens center toward the periphery according to a rule called the Fresnel zone, and this has multiple focal points by using different orders of diffracted light generated from that structure. In particular, when using a diffractive multi-focal lens as a contact lens or an intraocular lens, normally, 0th order diffracted light is the focal point for far vision, and +1 order diffracted light is the focal point for near vision. By distribution of this diffracted light, it is possible to make a bifocal lens having focal points for far and near vision. The general Fresnel zone constitution is basically the zone pitches having the zone radius determined by Equation 1 below. This Equation 1 is hereafter called a Fresnel zone setting equation. Besides, the zone radius and the zone diameter refer to the radius of the zone outer diameter.\n                              r          n                =                              nK            P                                              [                  Equation          ⁢                                          ⁢          1                ]            \nrn is the outer diameter radius of the nth zone obtained from Equation 1. K is a constant. P is addition power for setting the focus point of first order diffracted light with the focus point of 0th order diffracted light as a reference, and by varying this, it is possible to change the focal point position of the first order diffracted light.\nFor example when the focal point by 0th order diffracted light is a focal point for far vision, and first order diffracted light is set as the focal point for near vision, when P (the addition power noted above) is made larger, the focal point position for near vision moves closer to the lens. Specifically, when using that lens for the human eye, objects that are closer become visible. Conversely, when P is made smaller, the focal position for near vision recedes away from the lens. In this case, when the lens is used in the human eye, the near points that are visible recede away.\nFor patients with advanced presbyopia, or patients who have an intraocular lens inserted, power of accommodation of the crystalline lens decreases or is lost, so it is preferable to use a lens for which the focal point is matched in the nearer direction as with the former example. In other words, an item is needed for which the addition power is set to be large. On the other hand, for patients for which the power of accommodation has not decreased that much, even if the near focal point position is not made that near, it is possible to see near objects by joint use with one's own residual power of accommodation, so there are cases when large addition power does not need to be set. Taking into consideration the status of the eyes of these patients, it is possible to obtain bifocal lenses that can be suitably used at different required powers for each patient by setting P.\nFurthermore, in recent years, with the goal of improving visual performance in the intermediate region between two focal points of a focal point for far vision and a focal point for near vision, diffractive multi-focal lenses with a focal point set for intermediate vision have been proposed. This diffractive multi-focal lens with three or more focal points set has a plurality of types of reliefs for which their respective diffractive primary lights give mutually different focal point distances formed overlapping, having a synchronous structure for which the grating pitches for the reliefs overlap with each other periodically. As specific examples, for example, disclosed previously by this patent applicant, there are Japanese Unexamined Patent Publication No. JP-A-2010-158315 (Patent Document 1) and PCT Japanese Translation Patent Publication No. JP-A-2013-517822 showing the subordinate concepts thereof (Patent Document 2).\nHowever, with the diffractive multi-focal lens with three or more focal points set using the background art constitution, there was the problem that it was difficult to sufficiently ensure the degree of freedom for tuning the optical characteristics respectively requested for the plurality of focal points. In particular, with the inventions noted in Patent Documents 1 and 2, when setting the plurality of focal points, effective tuning technology was not yet established by which while adjusting the light intensity of each focal point, there is suppression of noise form peaks due to multi-order light or the like that is generated secondarily on the optical axis other than the target focal points.\nAlso, with the invention noted in Patent Document 1, when setting the intermediate focal point between the far focal point and the near focal point, there was the problem that it is difficult to set the intermediate focal point to an optional target position on the optical axis."}
{"text": "1. Field of the Invention\nThe invention relates to a semiconductor process, and more particularly to a fabrication process to form a multi-tip floating gate.\n2. Description of the Related Art\nMemory devices for non-volatile storage of information are currently in widespread use, in a myriad of applications. A few examples of non-volatile semiconductor memory include read only memory (ROM), programmable read only memory (PROM), erasable programmable read only memory (EPROM), electrically erasable programmable read only memory (EEPROM) and flash EEPROM.\nAn advantage of EPROM is that it is electrically programmed, but for erasing, requires exposure to ultraviolet (UV) light.\nIn many circuit designs it is desirable to have a non-volatile memory device that can be erased and reprogrammed in-circuit, without the need to remove the device.\nEEPROM devices have the advantage of electrical programming and erasing, achieved by charging and discharging actions controlled by the control gate. The actions also affect the conductivity of the channel between source and drain.\nOne of the advantages of flash memory is its capacity for block-by-block memory erasure. Furthermore, memory erasure is fast, normally taking just 1 to 2 seconds for the complete removal of a whole block of memory. Another advantage of flash memory is its low power consumption. The voltages of a control gate, a source, and a drain are adjusted to program or erase in a split gate flash memory.\nFIGS. 1a to 1c are cross-sections of the conventional method for fabricating a floating gate of a split gate flash memory.\nIn FIG. 1a, a silicon substrate 101 is provided. A gate oxide layer 102, a doped polysilicon layer 103, and a nitride layer 104 having an opening 105 are sequentially formed on the silicon substrate 101.\nIn FIG. 1b, the doped polysilicon layer 105 exposed by the opening 105 is oxidized to form an oxide layer 106, and an edge of the oxide layer 106 is a Bird\"\"s Beak shape edge.\nIn FIG. 1c, the nitride layer 104 is removed. The doped polysilicon layer 103 is anisotropically etched to form a floating gate 103a using the oxide layer 106 as an etching mask.\nA split gate flash memory is completed after a control gate is formed on the floating gate and the silicon substrate 101 is implanted to form source/drain devices.\nIn the program step, high voltage is applied between the source and control gate. The high voltage applied to the source goes to the floating gate by the electric capacity coupling, and a high electric field is produced on the film gate oxide layer. The voltage is injected into the floating gate through the film gate oxide layer from the drain.\nIn the erase step, high voltage is applied between the drain and the control gate. A high electric field is produced on the film gate oxide layer by the electric capacity coupling. The voltage is injected into the drain through the film gate oxide layer from the floating gate, such that the gate oxide layer is damaged by the high voltage.\nWhen the edge of the floating gate is a tip, the electrical field is easily concentrated, and the point is easily discharged. If the point discharge is increased, the erasing effect is strong.\nIn addition, the die size is larger due to the addition of programming circuitry and there are more processing and testing steps involved in the manufacture of these types of memory devices.\nThe present invention is directed to a method for fabricating a multi-tip floating gate to increase the erasing effect.\nAccordingly, the present invention provides a method for fabricating a floating gate. A semiconductor substrate is provided, on which a gate dielectric layer, a conducting layer, a first insulating layer, and a patterned hard mask layer with an opening are sequentially formed, such that the opening exposes the first insulating layer. The first insulating layer and the conducting layer are sequentially etched to form a round-cornered trench using the patterned hard mask layer as a mask. The patterned hard mask layer is removed. A second insulating layer is formed in the round-cornered trench. The first insulating layer and the exposed conducting layer are sequentially removed using the second insulating layer as a mask, and the conducting layer covered by the second insulating layer remains as a floating gate.\nAccordingly, the present invention also provides a method for fabricating a floating gate. A semiconductor substrate is provided, on which a gate dielectric layer, a conducting layer, a first insulating layer, and a patterned hard mask layer with an opening are sequentially formed, such that the opening exposes the first insulating layer. The first insulating layer and the conducting layer are sequentially anisotropically etched to form a bottom corner trench using the patterned hard mask layer as a mask, wherein the round-cornered trench does not expose the semiconductor substrate. The patterned hard mask layer is removed. A second insulating layer is formed on the first insulating layer, and the round-cornered trench is filled with the second insulating layer, wherein the material of the second insulating layer is different from that of the first insulating layer. The second insulating layer is planarized until the first insulating layer is exposed, and the second insulating layer remains in the round-cornered trench. The first insulating layer and the exposed conducting layer are sequentially anisotropically etched until the gate dielectric layer is exposed, and the conducting layer covered by the second insulating layer remains as a floating gate."}
{"text": "It is highly desirable for tires to have good wet braking performance and good braking performance on snow and ice. It has traditionally been very difficult to improve simultaneously a tire's braking characteristics on wet surfaces and snow and ice surfaces. This is because a good wet braking tread formulation needs to provide a high energy dissipation at low temperatures between about −20° C. and about 0° C. at a frequency of 10 Hz. Such a high energy dissipation is related to a stiffening of the rubber tread at temperatures between −10° C. and −40° C. at 10 Hz, which negatively impacts braking on snow and ice.\nIt would therefore be desirable to have a rubber composition for tires wherein the wet braking performance and braking on ice and snow are simultaneously improved."}
{"text": "Integrated circuits are made possible by processes which produce intricately patterned material layers on substrate surfaces. Producing patterned material on a substrate requires controlled methods for removal of exposed material. Chemical etching is used for a variety of purposes including transferring a pattern in photoresist into underlying layers, thinning layers, or thinning lateral dimensions of features already present on the surface. Often it is desirable to have an etch process that etches one material faster than another facilitating, for example, a pattern transfer process. Such an etch process is said to be selective to the first material. As a result of the diversity of materials, circuits, and processes, etch processes have been developed with a selectivity towards a variety of materials.\nEtch processes may be termed wet or dry based on the materials used in the process. A wet HF etch preferentially removes silicon oxide over other dielectrics and materials. However, wet processes may have difficulty penetrating some constrained trenches and also may sometimes deform the remaining material. Dry etches produced in local plasmas formed within the substrate processing region can penetrate more constrained trenches and exhibit less deformation of delicate remaining structures. However, local plasmas may damage the substrate through the production of electric arcs as they discharge.\nThus, there is a need for improved systems and methods that can be used to produce high quality devices and structures. These and other needs are addressed by the present technology."}
{"text": "Semiconductor devices are used in a variety of electronic applications, such as personal computers, cell phones, digital cameras, and other electronic equipment, as examples. Semiconductor devices are typically fabricated by sequentially depositing insulating or dielectric layers, conductive layers, and semiconductive layers of material over a semiconductor substrate, and patterning the various material layers using lithography to form circuit components and elements thereon.\nThe semiconductor industry continues to improve the integration density of various electronic components (e.g., transistors, diodes, resistors, capacitors, etc.) by continual reductions in minimum feature size, which allow more components to be integrated into a given area. These smaller electronic components also require smaller packages that utilize less area than packages of the past, in some applications.\nOne smaller type of packaging for semiconductors is a flip chip (FC) ball grill array (BGA) package, in which semiconductor die are placed upside-down on a substrate and bonded to the substrate using micro-bumps. The substrate has wiring routed to connect the micro-bumps on the die to contact pads on the substrate that have a larger footprint. An array of solder balls is formed on the opposite side of the substrate and is used to electrically connect the packaged die to an end application.\nCorresponding numerals and symbols in the different figures generally refer to corresponding parts unless otherwise indicated. The figures are drawn to clearly illustrate the relevant aspects of the embodiments and are not necessarily drawn to scale."}
{"text": "The present invention relates to pumps in general, and more particularly to plunger-type pumps.\nThere are already known and in widespread use various constructions of pumps which, depending on the requirements of such pumps, such as the pressure conditions, types of media to be pumped, volume to be pumped, and whether the operation is to be intermittent or continuous, have different shapes and operating parameters. Among such pumps, there has been already proposed a plunger-type pump in which a plunger is mounted in a pump housing and extends into a pumping chamber defined by the latter for reciprocation therein. In a pump of this type, one-way valves are to be provided both at the inlet side and on the outlet side of the pump and, in this connection, it has already been proposed to construct the inlet valve as a valve sleeve which is mounted on the plunger for movement longitudinally thereof and for limited reciprocation therewith. Then, the sleeve valve subdivides the pumping chamber into a suction space and a pumping space, the former surrounding, and the latter being located within the sleeve valve.\nPumps of the latter type achieve excellent results while being simple. However, experience with this type of pumps has shown that the seaing of the interface between the sleeve valve and the plunger leaves much to be desired. Usually, such a seal is constructed as an interface seal and is achieved by selecting the proper tolerance between the inner surface of the valve sleeve and the outer surface of the plunger, and by properly selecting the materials of the two components which bound the interface. While, under many circumstances, this type of seal may work to satisfaction, it has been established that this type of seal cannot be used with certain media to be pumped, particularly with aggressive media. A further limitation to the use of this type of seal resides in the fact that this type of seal can only be used up to a certain pressure existing in the pumping chamber. While it is not absolutely impossible to properly select the tolerances between the abovementioned surfaces for any pressure region or pressure differential existing across the sleeve valve, it has been found that the high precision which is required for high pressure involves costs which are prohibitive."}
{"text": "1. Field of the Invention\nThe present invention relates to a ventilator for cleaning the air inside the rooms of buildings.\n2. Prior Art\nPresently, there are many different types of ventilators for circulating the air between the inside and the outside of the rooms of a building. One type of the ventilators uses a plurality of pivoting slats. These slats are arranged in cross formation and installed in vertical and horizontal directions in a square frame. Both ends of each slat are attached to the opposite or facing side edges of the square frame in a pivotal fashion.\nIn this type of conventional ventilator, both ends of each slat are kept in contact perpendicularly with the inner wall surfaces of the side edges of the square frame. Accordingly, projecting pieces are used for mounting the slats. The projecting pieces are provided so as to project perpendicularly from the inside wall surfaces of the frame, and holes formed at both ends of the slat are fitted over these projecting pieces so that the slats can pivot.\nHowever, if the frame is round in shape and not a square, it is difficult to install a plurality of slats in a pivotal fashion and in a cross formation inside the round frame."}
{"text": "Many of the answers that mankind and scientists are searching for lay deep underground in various layers that have formed over many millions of years. These inaccessible areas continue to hide the answers and proofs that many scientists and researchers seek.\nIn addition, in support of exploration of the planetary bodies in our solar system, it will be necessary to install a measurement unit, such as a lunar seismometer for measuring moonquakes, within the lunar surface. In order to isolate the seismometer from solar radiation and the extremes of heat and cold, it will be necessary to sink the seismometer into the regolith without leaving a hole above the unit for intrusion from the surface of the complicating ambient conditions.\nIn support of research on deep geologic formations, the way that scientists can sometimes try to access a small amount of this otherwise inaccessible material for analyses is to drill and bring up a core sample of the material. In some locations, core sampling is considered impractical or even destructive of the condition sought to be measured. Even if scientists can obtain a small core sample, their act of removing this material from its underground location will often change the material and contaminate the answers being sought.\nIt is a known uncertainty of physics and science that the act of retrieving and moving a material to observe it could change the properties of the material. By bringing a core sample out of the earth to a laboratory, the sample is subject to changes in ambient pressure, temperature, moisture content, and structure. The core sample is now in an altered state due to a changed relationship with other ground materials, the level of magnetism and other earth energies that pass through it, the exact directional position in the ground, the exposure to sunlight and artificial light, and a whole host of various electric magnetic pollution and signals and gases and toxins, etc."}
{"text": "The prior art discloses simple surveying appliances having a sighting device, with the aid of which a spatial point is sighted manually and the sighting direction is then altered manually to a next spatial point to be surveyed, for example by means of adjusting screws on a theodolite drive.\nDE 196 48 626 discloses a method and an apparatus for area surveying with a laser rangefinder having a laser transmitter and a laser receiver. The laser rangefinder is mounted on a stand. The apparatus furthermore comprises a tilting and rotating device for orientation and direction measurement, a telescopic sight and also an electronic evaluation unit for angle data capture, distance data capture and data transfer to a computer. For surveying a space, the appliance is positioned at a central location in the space, from where all spatial and/or area corner points to be detected can be sighted and impinged upon by the laser beam. In accordance with the disclosure of DE 196 48 626, the spatial points to be surveyed are in this case each sighted individually, if appropriate—in the case of relatively large distance—with the observation being supported by means of a telescopic sight.\nA similar apparatus and associated surveying method are disclosed in DE 44 43 413, the supplementary published patent application DE 195 45 589 and WO 96/18083, which claims the priority of DE 44 43 413. These describe a method and an apparatus for surveying and marking on distant lines, areas or in at least partly closed spaces. One or more relevant spatial points are surveyed according to in each case two solid angles and the distance relative to a reference location by means of a laser distance measuring appliance mounted in a cardan-type fashion. The laser distance measuring appliance can be swiveled about two mutually perpendicular axes equipped with goniometers. In accordance with one embodiment described in said documents, spatial points to be surveyed are headed for manually and marking points are calculated from the survey data, on the basis of a prescribed relative relation between surveying and marking, said marking points then being moved to automatically by the measuring and marking apparatus.\nKnown construction surveying appliances typically comprise a base, an upper part mounted so as to be able to rotate about an axis of rotation on the base, and a sighting unit, mounted so as to be able to swivel about a swivel axis, with a laser source, which is designed to emit a laser beam, and an imaging detector, for example equipped with an orientation indicating functionality for indicating an orientation of the sighting unit with respect to a spatial point as a sighting point, and also with a distance determining detector for providing a distance measuring functionality. By way of example, the orientation indicating functionality may be a reticle in the viewfinder of a camera as imaging detector.\nModern, automated construction surveying appliances furthermore comprise rotary drives, which make the upper part and/or the sighting unit drivable in a motorized manner, goniometers and, if appropriate, inclination sensors for determining the spatial orientation of the sighting unit, and also an evaluation and control unit, which is connected to the laser source, the distance determining detector and also the goniometers and, if appropriate, inclination sensors.\nIn this case, the evaluation and control unit is equipped, by way of example, with a display having input means for inputting control commands from a user on the display (e.g. touchscreen) or what is known as a joystick that is directable, for the purpose of altering the orientation of the sighting unit by directing the joystick, and for presenting an image from the imaging detector or the camera on the display, wherein the orientation of the sighting unit can be indicated by means of the orientation indicating functionality on the display, e.g. by means of overlaying. Functionalities are known in which the input means on the display are in the form of arrows, the marking and touching of which enable a user to alter the orientation of the sighting unit in a horizontal or vertical direction.\nComputer technology reveals remote control units that are equipped with motion sensors and the movement of which is converted into an alteration in the position of what is known as a cursor or indicator arrow on a computer screen, in the form of what is known as a “computer mouse” or in the form of controllers for computer games.\nThe document JP 2004 108 939 A describes a system for controlling a total station by moving a remote control. The remote control unit contains acceleration and gravity sensors that detect movements by the remote control. The detected movements are converted into control commands and sent to the total station. What is not described in this case is scalability of the sensitivity level, i.e. of the transmission ratio between a movement by the remote control unit and a resultant speed and/or extent of the change in the orientation of the total station. The lack of scalability is particularly disadvantageous for fine orientation to a target."}
{"text": "Efforts have been made to make a thin composite resin-made container body to protect the environment and reduce costs. As an example, disclosed in Japanese Unexamined Patent Application No. 2000-313443 is a composite container that supplements rigidity of the container body reduced by being thinner, using a cylindrical paper label, in other words, a paper cylinder that covers the trunk of the container body.\nConventionally, this kind of composite container is made by forming the container body, then covering the container body with a separately made paper cylinder; it can take hours to combine the two so a limit has come to cost reductions.\nHowever, as disclosed in Japanese Unexamined Patent Application Publication No. 2011-116082, it is common for this kind of container body to be made by moving a preform P in a lateral direction (horizontal direction) to a two-piece, split-mold blow molding die BM that opens and closes by rocking, and blow molding after closing the mold."}
{"text": "The present invention relates to a personal care apparatus and, in particular, to an apparatus for collecting, containing and disposing of waste matter, during the replacement of a stoma bag for those people having had surgical operations for removing a bladder or colon.\nSurgical procedures such as colostomies, cystostomies, urostomies and ileostomies involve rerouting of the colon or ureter so that waste materials can be discharged through an artificial opening formed in the ostomy patient's (also referred to herein as ostomate) body. This artificial opening, called a stoma, is typically located in the abdomen and may be about 0.5 to 1.0 inch or more in diameter. A stoma is connected directly to the kidney by a rerouted ureter or to the intestines by the rerouted colon and the discharge of waste matter in each instance is made through the stoma.\nThe new artificial opening, or stoma, made on the abdominal wall, has no voluntary sphincter control by the ostomate. Collection of involuntary waste seepage is by bag or other receptacle attached to the ostomate's body. Disposable and semi-disposable bags for such purposes are available from several commercial sources. The disposable bags come as assembled units which are applied to the body. Such bags are attached to the body over the stoma by means of a ring-like member known as a wafer or face plate and a sealing ring of double sided adhesive which must be assembled and positioned over the stoma with precision to achieve a reliable attachment and seal. It is essential that this ring be mounted directly and concentrically over the stoma and further, that it be tightly secured to the body, otherwise, leakage of waste onto the surrounding skin area can occur. After removal of the filled bag, the common practice in the past has been to first clean and dry the stoma and surrounding skin area before application of the sealing ring to the skin, followed by connection of the face plate to the sealing ring and then attachment of a new bag to the face plate. It is most important that the skin area surrounding the stoma be kept clean and dry during replacement of the bag to avoid infection and to achieve a good connection between the skin, sealing ring and face plate. Adhesive tape is sometimes applied over the edges of the face plate and sealing ring for extra reinforcement. Inadvertent disconnection of a bag can cause considerable inconvenience and embarrassment to an ostomate.\nIn the replacement of the ostomy appliance, the general procedure is first to remove the ring and the receptacle. The ostomate applies an absorbent material of choice on top of the stoma to absorb any waste matter and prevent it from leaking onto the surrounding skin prior to the application of the adhesive. The surrounding skin must be cleaned and dried. This is generally accomplished by the use of a commercially available “prep” agent, which removes the skin oils, cleans the skin, and leaves a dry skin surface. Upon application of the adhesive, the absorbent material is removed from the stoma and a wafer or ring is placed over the stoma in a precisely centered position. Generally, one must remain in a standing position and using a mirror to achieve the centered position. The ostomy appliance is then attached to the ring and tightly secured to the skin.\nSeveral patents and publications describe devices designed to assist ostomates in the replacement of an ostomy device, including the following patents incorporated herein by reference.\nU.S. Pat. No. 6,409,709 describes an ostomy-changing facilitating device for absorbing fluids from the stoma during the change of the more permanent ostomy appliance. The ostomy-changing facilitating device includes a tubular member with absorbent material filling the bore of the tubular member.\nU.S. Pat. No. 8,343,119 describes a cap device for home use to temporarily cover the stoma resulting from a urostomy, colostomy or ileostomy procedure. It provides a means of preventing leakage during routine cleaning around the stoma during the ostomy pouch replacement process.\nU.S. Pat. No. 4,187,850 describes a medical-surgical method and apparatus for facilitating the replacement of ostomy appliances. A hollow cylinder, filled with sterile absorbent material is used to center the sealing ring prior to the replacement of the ostomy appliance.\nU.S. Pat. No. 4,344,433 describes an appliance for facilitating replacement of a waste bag by ostomy patient comprising a base, an upright post removeably seated in said base, and a slideable collar supported on an abuttment on said post. One end of the post includes an opening which covers the stoma and collects any seepage during the replacement procedure. This particular device requires that the face plate and adhesive sealing ring be preassembled on the post and collar, after which the post is removed from the base and positioned over the stoma.\nThe aforementioned prior art do not include all the features and advantages of the present invention, more specifically, the present device provides an absorbent waste receptacle with a dual absorbency chamber designed to effectively collect, contain and dispose of both urine and excrement; whereas the prior art typically provides only a hollow chamber, or a chamber entirely packed with absorbent material. The device of the present invention does not require preassembly of an ostomy appliance (such as pouch and wafer) or preassembly of the device itself, therefore allowing an ostomate to take ample time to prepare and perform the pouch replacement process. Furthermore, the device of the present invention provides an absorbent waster receptacle that is customizable to a particular ostomate's stoma as well as being disposable for a sanitary and discreet disposal of waste."}
{"text": "1. Field of the Invention\nThe invention relates to a device for lighting a group of selectively positioned candles and more particularly to a means for lighting a group of birthday type candles that are threadably mounted on an elongated igniting fuse line or string.\n2. Description of the Prior Art\nBirthday candles are well known and have been for many years used as a means to light up birthday cakes. The number of candles that are used on a cake vary, of course, depending upon the age of the person whose birthday it might be. A birthday candle is generally formed having an elongated wax body 60 mm long with a diameter of approximately 5 mm and is provided with an internally positioned wick that extends outwardly from the top of the wax body.\nThese candles, when placed on the surface of the cake, are often arranged to form various decorative design configurations by placing them in straight or curved rows and circles. They can also be arranged to form, for example, heart shapes, stars or positioned to define letters, and so on.\nHowever, when a large group of candles are positioned on a cake they often present several problems when they are to be lit with a match. Usually one does not know which candle to light first so that he or she does not get burned in the process of lighting the remaining candles. Birthday type candles are provided with an extended wick formed from a cotton material which is very soft and sometimes bent to the side of the candle. This too creates a problem in lighting the wick. After several candles are lit it becomes increasingly harder to reach the remaining unlighted candles without getting burned. Further, when trying to light a large group of candles the first few lighted candles have already started to melt which causes the wax to drop on the icing of the cake. Often the flames of the first lighted candles become rather large.\nAccordingly, there is a need to provide a simple and safe device and/or method of lighting a group of closely mounted candles that are commonly found with respect to birthday cakes."}
{"text": "Human interleukin 12 (IL-12) has recently been characterized as a cytokine with a unique structure and pleiotropic effects (Kobayashi, et al. (1989) J. Exp Med. 170:827-845; Seder, et al. (1993) Proc. Natl. Acad. Sci. 90:10188-10192; Ling, et al. (1995) J. Exp Med. 154:116-127; Podlaski, et al. (1992) Arch. Biochem. Biophys. 294:230-237). IL-12 plays a critical role in the pathology associated with several diseases involving immune and inflammatory responses. A review of IL-12, its biological activities, and its role in disease can be found in Gately et al. (1998) Ann. Rev. Immunol. 16: 495-521.\nStructurally, IL-12 is a heterodimeric protein comprising a 35 kDa subunit (p35) and a 40 kDa subunit (p40) which are both linked together by a disulfide bridge (referred to as the “p70 subunit”). The heterodimeric protein is produced primarily by antigen-presenting cells such as monocytes, macrophages and dendritic cells. These cell types also secrete an excess of the p40 subunit relative to p70 subunit. The p40 and p35 subunits are genetically unrelated and neither has been reported to possess biological activity, although the p40 homodimer may function as an IL-12 antagonist.\nFunctionally, IL-12 plays a central role in regulating the balance between antigen specific T helper type (Th1) and type 2 (Th2) lymphocytes. The Th1 and Th2 cells govern the initiation and progression of autoimmune disorders, and IL-12 is critical in the regulation of Th1-lymphocyte differentiation and maturation. Cytokines released by the Th1 cells are inflammatory and include interferon γ (IFNγ), IL-2 and lymphotoxin (LT). Th2 cells secrete IL-4, IL-5, IL-6, IL-10 and IL-13 to facilitate humoral immunity, allergic reactions, and immunosuppression.\nConsistent with the preponderance of Th1 responses in autoimmune diseases and the proinflammatory activities of IFNγ, IL-12 may play a major role in the pathology associated with many autoimmune and inflammatory diseases such as rheumatoid arthritis (RA), multiple sclerosis (MS), and Crohn's disease.\nHuman patients with MS have demonstrated an increase in IL-12 expression as documented by p40 mRNA levels in acute MS plaques. (Windhagen et al., (1995) J. Exp. Med. 182: 1985-1996). In addition, ex vivo stimulation of antigen-presenting cells with CD40L-expressing T cells from MS patients resulted in increased IL-12 production compared with control T cells, consistent with the observation that CD40/CD40L interactions are potent inducers of IL-12.\nElevated levels of IL-12 p70 have been detected in the synovia of RA patients compared with healthy controls (Morita et al (1998) Arthritis and Rheumatism. 41: 306-314). Cytokine messenger ribonucleic acid (mRNA) expression profile in the RA synovia identified predominantly Th1 cytokines. (Bucht et al., (1996) Clin. Exp. Immunol. 103: 347-367). IL-12 also appears to play a critical role in the pathology associated with Crohn's disease (CD). Increased expression of INFγ and IL-12 has been observed in the intestinal mucosa of patients with this disease (Fais et al. (1994) J. Interferon Res. 14:235-238; Parronchi et al., (1997) Am. J. Path. 150:823-832; Monteleone et al., (1997) Gastroenterology. 112:1169-1178, and Berrebi et al., (1998) Am. J. Path 152:667-672). The cytokine secretion profile of T cells from the lamina propria of CD patients is characteristic of a predominantly Th1 response, including greatly elevated IFNγ levels (Fuss, et al., (1996) J. Immunol. 157:1261-1270). Moreover, colon tissue sections from CD patients show an abundance of IL-12 expressing macrophages and IFNγ expressing T cells (Parronchi et al (1997) Am. J. Path. 150:823-832).\nDue to the role of human IL-12 in a variety of human disorders, therapeutic strategies have been designed to inhibit or counteract IL-12 activity. In particular, antibodies that bind to, and neutralize, IL-12 have been sought as a means to inhibit IL-12 activity. Some of the earliest antibodies were murine monoclonal antibodies (mAbs), secreted by hybridomas prepared from lymphocytes of mice immunized with IL-12 (see e.g., World Patent Application Publication No. WO 97/15327 by Strober et al.; Neurath et al. (1995) J. Exp. Med. 182:1281-1290; Duchmann et al. (1996) J. Immunol. 26:934-938). These murine IL-12 antibodies are limited for their use in vivo due to problems associated with administration of mouse antibodies to humans, such as short serum half life, an inability to trigger certain human effector functions and elicitation of an unwanted immune response against the mouse antibody in a human (the “human anti-mouse antibody” (HAMA) reaction).\nIn general, attempts to overcome the problems associated with use of fully-murine antibodies in humans, have involved genetically engineering the antibodies to be more “human-like.” For example, chimeric antibodies, in which the variable regions of the antibody chains are murine-derived and the constant regions of the antibody chains are human-derived, have been prepared (Junghans, et al. (1990) Cancer Res. 50:1495-1502; Brown et al. (1991) Proc. Natl. Acad. Sci. 88:2663-2667; Kettleborough et al. (1991) Protein Engineering. 4:773-783). However, because these chimeric and humanized antibodies still retain some murine sequences, they still may elicit an unwanted immune reaction, the human anti-chimeric antibody (HACA) reaction, especially when administered for prolonged periods.\nA preferred IL-12 inhibitory agent to murine antibodies or derivatives thereof (e.g., chimeric or humanized antibodies) would be an entirely human anti-IL-12 antibody, since such an agent should not elicit the HAMA reaction, even if used for prolonged periods. However, such antibodies have not been described in the art and, therefore are still needed."}
{"text": "The present invention relates to electrical connecting arrangements in general, and more particularly to a connector for connecting two electric conductors with one another.\nThere are already known various constructions of conductor connecting arrangements of the above type, among them such which include at least one current bar having at least one end portion which is bent backwards upon itself, and at least one clamping body which can be slid onto the bent-back end portion of the current bar from a central region of the current bar and serves for connecting an electrical conductor to the current bar. A conductor connector arrangement of this type, which is of a particular interest for use in high-current terminals or clamps, is known, for instance, from the German published patent application DE-A-No. 34 08 006. In conductor connector arrangements of this type, it is often necessary to connect relatively thick-wire conductors, so that it is essential to be able to lay the conductor freely onto the current bar first, and only then to slide the clamping body over the current bar and the conductor.\nIt is further known in principle to slide the clamping body of a screw connector arrangement in a guided manner on the current bar and over the conductor positioned on the current bar to the clamping location. Herein, the clamping body is guided to a certain extent and more or less in the correct position on the current bar itself, as disclosed in the Swiss patent CH-PS No. 392 658, or guides for the clamping body are provided in a complicated manner by correspondingly constructing and configuring the electrically conductive parts themselves and providing special guiding elements and arresting holders.\nSuch guiding possibilities are not conventionally provided or available in conductor connectors of the type discussed above. Moreover, it is to be considered that the positioning of the clamping body from the longitudinally central region of the current bar onto the bent-back end portion of the current bar requires, to begin with, a considerable amount of skill and that it would be especially difficult to perform this operation if it was desired to associate such a conductor connector with or accommodate it in a sort of an insulating material housing, for instance, for reasons of prevention of contact with electric current carrying elements"}
{"text": "Systems and methods that provide map-based directions for mobile devices and other computers are well known. For example, Google Maps for Mobile allows mobile devices, such as cell phones and PDAs, to display maps and text-based directions in response to user requests.\nSome software, including Google Maps for Mobile, also allows users to display their progress relative to a map in real time. For example, if a mobile device is able to determine its geographic location, the device may download or otherwise access maps corresponding with its location and then display its location on the map. Mobile devices typically calculate their position via an internal or external GPS component, or by calculating their position relative to one or more cell phone towers.\nRelatively complete and accurate maps are available with respect to many parts of the world. However, in many other areas of the world, the map data is incomplete and sometimes even incorrect, particularly when it comes to undeveloped or rapidly-developing areas. For example, a map database may not reflect the fact that a road has been recently built, closed or rerouted. Moreover, the database may be aware that a road is present, but it may not store the street name. Some streets will not even have a name.\nOne method of coping with the problem of incomplete maps is for mapping companies to retain people who drive around with GPS devices, manually detect road changes and enter new roads into the database. In addition to being relatively expensive, such a method of uploading data to a map database has other disadvantages as well."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for driving a display made up of a current-driving type light-emitting device which displays various pieces of information, results of measurement, moving pictures, or still pictures and to a circuit employing the above method and to portable electronic devices incorporating the circuit and more particularly to the method for driving the display which is used as a display device of computers such as notebook computers, palm-sized computers, pocket computers, or a like, of portable electronic devices such as a PDA (Personal Digital Assistant), a portable cellular phone, or a PHS (Personal Handy-phone System) and to the circuit employing the above method and to the portable electronic devices being equipped with driving circuits for the display.\nThe present application claims priority of Japanese Patent Application No.2001-285838 filed on Sep. 19, 2001, which is hereby incorporated by reference.\n2. Description of the Related Art\nSome types of displays are made up of current-driving type light-emitting devices. The displays of this kind conventionally include a display made up of an EL (Electroluminescnece) device, a display made up of an LED (Light-Emitting Diode), a VFD (Vacuum Fluorescent Display) including, in particular, an FED (Field Emission Display) being one of types of the VFDs, a PDP (Plasma Display Panel), or a like. Hereinafter, this type of the display is called a “self-emissive type display”.\nGenerally, the self-emissive type display tends to draw current more than a voltage-driving type liquid crystal display. That is, in the voltage-driving type liquid crystal display, since its liquid cell is a capacitive load, an amount of current consumed is as little as several mA. However, since the self-emissive type display emits light for every pixel and therefore consumes current every time it emits light, an amount of current consumed reaches 200 mA or more when high current volumes are consumed, for example, when an image is displayed at a high value of luminance. Therefore, when the self-emissive type display is used in a displaying section of a portable electronic devices to which power is supplied by a battery, dry cell, or a like, to keep operating time as long as possible, an amount of current consumed has to be reduced to a minimum. The portable electronic devices include notebook-type, palm-type, or pocket-type computers, PDAs, or portable cellular telephones, PHSs, or a like.\nThe portable cellular phone or the PHS has a waiting mode in which, though power is provided, a user does not perform any operation while waiting for an incoming call. The displaying section provides a waiting screen corresponding to a waiting mode.\nNot only the portable cellular phone or the PHS but also other portable electronic devices when, though power is provided, and a specified time has elapsed without any operations performed by a user, as shown in FIG. 8, is put into a screen save mode in which a specified character or diagram is displayed randomly on each portion on the display at every specified time and a moving object is displayed on the display. FIG. 8 shows that a current time indicated by the self-emissive type display in the portable cellular phone or the PHS is displayed in the screen save mode in a manner that a place for displaying time is changed at every specified time. The screen save mode is used to prevent occurrence of a phenomenon called an “image burn-in” state in which same characters or same diagrams continue to be displayed for a long time, even if power is turned OFF, a trace of the character or the diagram is left. The user, when such portable electronic devices are in the waiting mode or in the screen save mode, does not view the screen of a display carefully.\nHowever, conventionally, even when the user drives a screen of a display in a waiting mode or lowers luminance of each pixel in a screen save mode, the user employs the same driving method as is used for driving the ordinary screen on which the user views the screen carefully. For this reason, in the former case, power is drawn wastefully. Moreover, in the latter case, the display becomes dark as a whole and are hard to view and, therefore, if the user views the display unexpectedly, the user cannot confirm contents of the display immediately and misunderstands, in some cases, that power has not been provided.\nIn some cases, as shown in FIG. 9, the display screen of the conventional portable cellular phone or the PHS is made up of, for example, an upper display portion 1, a central display portion 2, and a lower display portion 3. In the upper display portion 1, a battery mark 1a indicating a charging state of a battery, an antenna mark 1b indicating whether or not the portable cellular phone or the PHS being presently used is in a service area of a wireless telephone system of a mobile communication network, or a like are displayed. In the central display portion 2, sentences of electronic mail, images being attached to the electronic mail, and images indicating contents being provided from various contents providers in a WWW (World Wide Web) server or a like are displayed. FIG. 9 shows an example in which map data is displayed in the central display portion 2. In the lower display portion 3, a menu key used to select a menu is displayed. Then, generally, in the central display portion 2 is displayed a detailed image, while in the upper display portion 1 and in the lower display portion 3 are displayed a character and/or a mark in a simplified manner. This is because, even if the character and/or mark is simplified, information can be fully transferred to the user of the portable cellular phone or the PHS.\nHowever, conventionally, even in the case of the upper display portion 1 and in the lower display portion 3 in which such simplified characters or marks are displayed, a same driving method as is used to display in the central display portion 2 in which detailed images are displayed is employed. This causes wasteful power consumption. The same inconveniences as described above occur in other portable electronic devices to which power is supplied by the battery or the dry cell such as the notebook-type, the palm-type, and pocket-type computers, the PDA or the like, though contents and portions displayed therein (for example, in a window), or a like are different."}
{"text": "Typically, polymers are used in the manufacture of printed wiring boards (PWB), by coating the PWB with a polymer coating. For solder masks the coating pattern covers conductor wiring traces at all locations except where electrical contact is to be made with the traces. The curing traces may also be printed on a metallic layer such as copper on an insulating substrate by means of polymer printing patterns defining areas which are to be etched in one primary imaging method. In other primary imaging methods, polymer printing patterns may define zones for accumulation of conductive materials forming conductor traces. Other printed objects are similarly manufactured by use of the polymer printing patterns.\nIn common use are three types of polymers, namely, (1) thermal curing epoxy, (2) ultra-violet (UV)-curing polymers generally classified as liquid polymers which are hardened by radiation, and (3) dry film polymers which are changed in solubility by radiation. The epoxies and UV-curables are applied as liquid coatings for curing in place as contrasted with the dry films.\nSome of the patterns formed by these polymers are formed mechanically such as by silk screening. Others are formed by photo-imaging techniques, where the photo response characteristics of certain polymers such as the UV curing liquid polymers are utilized.\nDry film solder mask photopolymers are supplied by two manufacturers, in roll form, consisting of an inner layer of photopolymer sandwiched between a carrier film of clear polyester and a polyolefin liner. They have problems of adhering to copper, and are not conveniently used for overlaying rough surfaces. Thus, for a solder mask coating on a typical PWB, wherein the metal conductors extend 0.003 to 0.004 inch above the base laminate, the dry film needs be conformed without intermediate air bubbles. The use of a roller laminator is not usually satisfactory, as air is trapped between the photopolymer and the PWB laminate, particularly between closely-spaced conductors. A vacuum chamber laminator is normally used in a plasticizing procedure to prevent air entrapment. The lamination cycle is as follows: A PWB and a section of dry film solder mask is inserted into the heated chamber; the chamber is evacuated, and when up to temperature the photopolymer is forced into contact with the PWB, effecting an air-free lamination. Lamination occurs at a temperature of the order of 200 degrees F., at which temperature the dry film solder mask photopolymer becomes tacky and adheres to the PWB surface. The solder mask pattern is attained by exposing the photopolymer to a strong UV light source through a photographic film, wherein the light hardens the exposed photopolymer. The polyester carrier film is then peeled away and the unhardened photopolymer is washed out in a solvent spray bath.\nThere are several shortcomings with the dry films listed below:\n1. The single photopolymer layer is expensive being of the order of four times the cost of liquid photopolymers; the equipment required to laminate is overly complicated and expensive; the labor required is excessive, as each processing step of laminating, exposing, and developing is overly lengthy.\n2. The laminating step forces photopolymer into circuit holes, and the like, with the result that with small holes of the order of 0.025 inch diameter, the photopolymer does not wash out, leaving the holes plugged.\n3. The laminating step forces photopolymer into large tooling holes and slots, leaving a puckered, striated coating of photopolymer, for the polyester carrier film is non-conforming and wrinkles around larger holes and sharp corners.\n4. The dry photopolymer is characterized by a lack of adhesion to metal conductors. For PWB having bare copper conductors with dry film solder mask thereover, the solder coating and hot-air leveling step is usually not satisfactory without copper surface pretreatment as with a black oxide coating; otherwise the solder mask separates from the conductors.\n5. The dry photopolymers need be applied in thick layers of standard thicknesses, and the thicker layers are not only expensive but also inhibit high resolution photo reproduction and require higher temperatures to photodevelop.\n6. The photopolymer is temperature sensitive, delaying the application sequence. After laminating at 200 degrees F., the photopolymer must cool to room temperature prior to exposure. After exposure, which raises the temperature again, the photopolymer must cool to room temperature before washing out unexposed photopolymer.\nSimilarly epoxy coatings have disadvantages since the initial liquid layers required need be cured in situ, generally at very high temperatures which can damage the printed wiring boards or other substrates being printed upon.\nThe liquid polymers provide problems in initial application since they do not tend to stay in place, and have conventionally been applied in very thin coatings, unsuitable for pinhole free coatings, and for covering the rough surfaces of printed wiring traces.\nThe different polymers each have advantageous features, but it is not known in the prior art how to overcome this disadvantage in producing printed patterns on substrates therewith.\nTypical of prior art polymer printing techniques used in PWB are those of U.S. Pat. No. 3,629,036--C. Isaacson, Dec. 21, 1971, where a microthin adhesive layer including a photopolymer solvent is interspersed between smooth copper surfaces and a dry film photoresist layer, which is thereby glued in place before photoimaging and developing the desired pattern.\nReference is made to U.S. Pat. No. 3,824,104 in which Kloczewski teaches a method for photoimaging a liquid photopolymer, wherein the image bearing photomask is separated from the photopolymer by a distance of eight mils during exposure, leading to a stated loss of resolution.\nReference is made to U.S. Pat. No. 4,260,675 in which Sullivan describes a method for photoimaging liquid photopolymer using a glass plate photomask with raised opaque pillars in contact only with portions of the PWB which will be free of hardened solder mask.\nEach of these patents and current practices in the art of PWB manufacture has characteristics which this invention seeks to improve.\nOne objective is to improve the adhesion of dry film to metal conductors.\nA second objective is a process to laminate a dry film to a highly irregular surface such as over PWB traces without use of a vacuum laminator.\nAnother objective is to reduce the time required to process dry film by negating the requirement for laminating at an elevated temperature.\nAnother objective is to achieve a photoimaged solder mask having the electrical and environmental characteristics of dry film, but at reduced cost.\nAnother objective is to achieve a photoimaged PWB solder mask in which UV-curable liquids are exposed with a photomask in contact with the liquid, and in which the solder mask coating is not thinned out over the metal conductors."}
{"text": "The present invention relates generally to strap retainers and, more particularly, to strap retainers for retention of shopping cart safety straps.\nIt is well known in the prior art to secure a nylon safety strap child restraint to a grocery shop cart in order to assure the safety of a child riding in the cart. It is important that such devices are securely attached to the shopping cart and it is desirable that the retainer does not become accidentally detached. In addition, the strap retainer should be resistant to damage, or removal, as a result of vandalism. This latter consideration is important because, in many cases, shopping carts are left in open parking lots for substantial periods of time thereby affording opportunities for vandalization.\nSeveral conventional devices have been used to secure a child safety strap to a grocery cart. In general, these devices are complicated, sometimes made of metal and subject to rust, and are relatively expensive to produce. An example of a complicated device is disclosed in U.S. Pat. No. 5,669,118 in which a complicated combination of retainer elements is utilized. This device, as well as some other conventional devices, is relatively difficult to use.\nAccordingly, there is a need for an efficient, low cost and effective technique for attaching a safety strap to a shopping cart.\nThe environment in which the device is to be used typically dictates the desirable characteristics of a suitable safety strap retainer. In a typical case, the safety strap retainer is fixed to a convenient shopping cart rail and is thereafter subject to somewhat rough usage. It is not unusual for shopping cart to be left outside, in wet weather, for long periods of time and a metal strap retainer tends to rust under these circumstances. Clearly, it would be desirable to have a strap retainer that is relatively weather resistant and able to withstand rough usage.\nAs another consideration, retail store employees often attach the strap retainer to the shopping cart. In this regard, it would be desirable if the retainer could be installed conveniently, without any need for a tool. In addition, a suitable strap retainer should be simple in construction and inexpensive to produce. Further, the strap retainer should be difficult to remove, after it has been installed, so that the retainer would not be susceptible to damage by vandals.\nFrom the foregoing it will be apparent that there is a need for a strap retainer for use in securing a safety strap to a shopping cart that is mechanically simple, easy to use, low in cost and resistant to easy removal."}
{"text": "1. Field of the Invention\nThe present invention relates to a color matching method and an image capturing device using the color matching method or an electronic apparatus provided with such image capturing device. More particularly, the present invention relates to a color matching method for handling an electronic digital image and an image capturing device using the color matching method or an electronic apparatus provided with such image capturing device.\n2. Description of the Related Art\nRecently, electronic apparatuses provided with a color image capturing device are widely used, for example, a digital camera, an electronic type video recorder, or a mobile phone. Generally, a semiconductor image sensor is provided in these electronic apparatuses for converting a captured image into an electrical signal. Due to the inherent limitations of the image sensor with respect to the color sensing performance, the color of an image outputted under such limitations is usually required to be adjusted. In addition, the color of the image is different depending on the light illumination conditions. For example, the perception of color for the human visual system is quite different in the moonlight and in the daylight outdoors. Therefore, it is important that the obtained electronic image is required to meet the perception of the human's visual sense in color processing.\nIn these electronic apparatuses, an image signal outputted from the image sensor will be subjected to a series of signal processing procedures in which the color of the obtained image will be corrected by performing a color matching process so as to output an image with a natural color.\nRegarding the color correction process, a number of conventional techniques that adopt a correction matrix to correct the color of an image are known. For example, U.S. Pat. No. 7,053,935 has disclosed one advantage over current electronic capture systems that the accuracy of color capture, which is a result of the match between the human visual system and the unique set of spectral response curves and the correction matrix used in the invention. However, the color correction result cannot be dynamically adjusted since the matrix is predefined, and thus a quick, accurate color correction result is difficult to be achieved, and a high-quality image cannot be offered.\nAs a result, in the conventional technique described above, since a fixed correction matrix is used to correct the color of the image, the color difference is large and cannot be rapidly converged to a predetermined value.\nTherefore, a color matching method capable of dynamically adjusting the color difference and rapidly converging the color difference to the predetermined value is required for achieving a color demonstration with high color accuracy."}
{"text": "A mobile robot is a machine which can work automatically. The mobile robot can be operated under the command of human operators or in pre-programmed programs, and can act according to principles set out by the artificial intelligence technology as well. This type of mobile robot can be used indoors or outdoors, and can be used in industry, business or household. For example, the mobile robot can be used to replace security guards to perform patrol, or replace people to clean the surface. The mobile robot can also be used to accompany family members or assist in doing office work.\nWhen mobile robot like cleaning robot, accompanying robot or guest-guiding robot moves under an operating mode, due to the complexity of the operating environments, the mobile robot may often hit obstacles so as to damage furniture and itself in the operating process because of an unfavorable obstacle avoidance performance thereof, thereby influencing normal operation of the mobile robot. Wheels of the mobile robot often get entangled with some flexible obstacles (e.g., cables, ropes, ribbons, leftover of cloth), so that the mobile robot cannot move, and even the mobile robot may tip over to cause safety accidents. For a cleaning robot, the cleaning system of the cleaning robot may be more likely to be entangled by flexible obstacles, for example, cleaning brushes may be entangled by flexible obstacles, or a dust collection component may be entangled or blocked by flexible obstacles.\nGenerally speaking, in order to detect obstacles on the floor such as ground in time, obstacle detection technologies often adopted on the existing mobile robot are mainly as follows:\nMechanical collision detection: a mechanical baffle connected with an electronic switch is arranged at the front of the bottom of the mobile robot, when the mobile robot comes into contact with obstacles, the electronic switch is converted into a connected state from a disconnected state, so as obstacles in the front can be detected. In this way, detection of obstacle can be performed only after collision occurs, and thus contribute to a poor user experience. Moreover, flexible obstacles are easily pushed away by the mobile robot, the electronic switch cannot be triggered because the collision force between robot and obstacles is too small, thereby leading to a missing detection.\nInfrared distance measurement detection: one or more infrared distance measurement sensors are arranged on the mobile robot. And obstacles in the front can be detected in the case of the detected distance is smaller than a preset threshold. However, infrared detection is greatly influenced by ambient lighting, and there is a greater dead zone within a short distance. Obstacles made from glass, light-absorbing or all-black materials are easily missed in detection, and the consistency of precisions is poor. Because of the structure limitations, the infrared sensor cannot be arranged on the lower part of housing, thus some obstacles which are lower than the infrared rays, e.g., flexible obstacles, may be missed in detection.\nUltrasonic distance measurement detection: one or more ultrasonic distance measurement sensors are arranged on the mobile robot. And obstacles in the front can be detected in the case of the detected distance is smaller than a preset threshold. However, on the one hand, ultrasonic waves are easily influenced by factors such as ambient temperature, material of reflectors and multipath propagation of sound waves. On the other hand, Because of the structure limitations, the ultrasonic sensor cannot be arranged on the lower part of housing, thus some obstacles which are lower than the ultrasonic waves, e.g., flexible obstacles, may be missed in detection. In addition, another limitation of ultrasonic measurement is that it depends on the effective reflection area of the obstacle, such that object with a size greater than a certain threshold can be detected, smaller object such as wires can be easily missed.\nTherefore, it can be seen that the detection of flexible obstacles is the problem that needs to be solved urgently for the existing mobile robot."}
{"text": "Recently, as being accompanied by the developments of portable electronic devices such as cellular phones and notebook-size personal computers, or as being accompanied by electric automobiles being put into practical use, and the like, small-sized, lightweight and high-capacity secondary batteries have been required. At present, as for high-capacity secondary batteries meeting these demands, non-aqueous secondary batteries have been commercialized, non-aqueous secondary batteries in which lithium cobaltate (e.g., LiCoO2) and the carbon-system materials are used as the positive-electrode material and negative-electrode material, respectively. Since such a non-aqueous secondary battery exhibits a high energy density, and since it is possible to intend to make it downsize and lightweight, its employment as a power source has been attracting attention in a wide variety of fields. However, since LiCoO2 is produced with use of Co, one of rare metals, as the raw material, it has been expected that its scarcity as the resource would grow worse from now on. In addition, since Co is expensive, and since its price fluctuates greatly, it has been desired to develop positive-electrode materials that are inexpensive as well as whose supply is stable.\nHence, it has been regarded promising to employ lithium-manganese-oxide-system composite oxides whose constituent elements are inexpensive in terms of the prices as well as which include stably-supplied manganese (Mn) in their essential compositions. Among them, a substance, namely, Li2MnO3 that includes tetravalent manganese ions alone but does not include any trivalent manganese ions making a cause of the manganese elution upon charging and discharging, has been attracting attention. Although it has been believed so far that it is impossible to charge and discharge Li2MnO3, it has come to find out that it is possible to charge and discharge it by means of charging it up to 4.8 V, according to recent studies. However, it is needed to further improve Li2MnO3 with regard to the charging/discharging characteristics.\nIn order to improve the charging/discharging characteristics, it has been done actively to develop xLi2MnO3.(1−x)LiMeO2 (where 0<“x”≦1), one of solid solutions between Li2MnO3 and LiMeO2 (where “Me” is a transition metal element). Note that it is feasible to write and express Li2MnO3 by a general formula, Li(Li0.33Mn0.67)O2, as well, and that it is said to belong to the same crystal structure (i.e., a layered rock-salt structure) as that of LiMO2. Consequently, there arises a case where xLi2MnO3.(1−x)LiMeO2 is set forth as Li1.33−yMn0.67−zMey+zO2 (where 0<“y”<0.33, and 0<“z”<0.67), too, but even any of the two methods for writing it down specify a composite oxide that possesses the same sort of crystal structure.\nFor example, Patent Literature No. 1 discloses a production process for solid solution between LiMO2 and Li2NO3 (where “M” is one or more kinds that are selected from Mn, Ni, Co and Fe, and “N” is one or more kinds that are selected from Mn, Zr and Ti). This solid solution is obtainable as follows: ammonia water is dropped to a mixed solution, in which salts of respective metallic elements that correspond to “M” and “N” are dissolved, until the pH becomes 7; an Na2CO3 solution is further dropped to it in order to deposit “M”-“N”-system composite carbonates; and the resulting “M”-“N”-system composite carbonates are calcined after mixing them with LiOH.H20.\nHowever, upon employing a secondary battery including Li2MnO3 as the positive-electrode active material, it is needed to activate the positive-electrode active material at the time of the first-time charging. Since the activation is accompanied by a large irreversible capacity, ions having moved to the counter electrode do not come back, and so there is such a problem that charging/discharging balance between the positive electrode and the negative electrode becomes imbalanced. With regard to the mechanism of this activation and to an obtainable capacity by means of the activation, it is the present situation that they have not been clearly clarified yet (see Non-patent Literature No. 1).\nMoreover, in a case where a particle diameter of Li2MnO3 is large, since only the particles' superficial layer is activated, it is believed that it is necessary to make the particle diameter of Li2MnO3 smaller in order to turn Li2MnO3 to be employed into an active material serving as battery in the total amount virtually. In other words, it is needed to develop convenient processes for synthesizing fine particles as well. For example, in Patent Literature No. 2, a process for synthesizing nano-order oxide particles is disclosed. In Example No. 3 of Patent Literature No. 2, MnO2 and Li2O2 are added to and are then mixed with a mixture, in which LiOH.H2O and LiNO3 are mixed in a molar ratio of 1:1; and they are turned into 300° C. molten salt after letting the mixture go through a drying step, thereby synthesizing lithium manganate (e.g., LiMn2O4) with a spinel structure, whose manganese has an average oxidation number that is equal to a valence number of 3.5."}
{"text": "Conventional computing devices typically include one to many conventional types of input/output (I/O) ports for communicating with connectable external devices such as mice, keyboards, wireless modems, thumb drives, hard drives, etc., as well as internal devices such as hard drives and SSD drives. Conventional computing devices typically further include subsystems for inputting and outputting audio and video streams such as music, videos, video chat and conferences, presentations, etc.\nHowever, the specifications for these I/O and multimedia interfaces and subsystems typically do not provide for security functions such as authentication and verification. Meanwhile, there are a number of applications such as corporate video conferencing that would greatly benefit from efficient provision and management of security over such multimedia interfaces and subsystems."}
{"text": "Semiconductor devices are used in a variety of electronic applications, such as personal computers, cell phones, digital cameras, and other electronic equipment, as examples. As technology has progressed, the demand for smaller semiconductor devices with improved performance has increased. As feature densities increase, the widths of the conductive lines, and the spacing between the conductive lines of back-end of line (BEOL) interconnect structures also need to scale smaller.\nA move is being made away from the traditional materials used in the past in semiconductor device designs, in order to meet these demands. To reduce the RC time delay, low dielectric constant (low-k) materials are being used as insulating materials, and there is a switch being made to the use of copper for interconnect materials, rather than aluminum. Advantages of using copper for semiconductor device interconnects include abilities to operate faster and manufacture thinner conductive lines because copper has lower resistivity and increased electromigration resistance compared to aluminum. Combining copper interconnects with low-k dielectric materials increases interconnect speed by reducing the RC time delay, for example.\nCopper interconnects are often formed using damascene processes rather than by direct etching. Damascene processes are typically either single or dual damascene, which includes forming openings by patterning and etching inter-metal dielectric (IMD) layers and filling the openings with copper. Because copper diffuses easily into some dielectric materials, especially some types of low-k dielectric materials, a diffusion barrier layer is usually deposited on the inner walls of the damascene opening before the copper is formed. Refractory metals such as tantalum (Ta) or titanium (Ti), or nitride compounds of these metals are used as materials of the diffusion barrier layer. However, there are some challenges in using refractory metals in the copper damascene structure since these metallic films have high resistance, thereby causing increased resistance in the copper lines and increased RC delay, especially in small, narrow features.\nAs the shrinkage of copper lines has progressed in recent years, there is a trend towards thinner films being used for the diffusion barrier layer. Physical vapor deposition (PVD) process used for depositing a thinner TaN/Ta barrier layer encounters difficulties in advanced scale of interconnection. Atomic layer deposition (ALD) process is the candidate to deposit a very thin diffusion barrier layer with uniform coverage, but the ALD method is disadvantageous in extremely low deposition rate and poor throughput. In addition, in manufacturing the TaN/Ta film, a problem occurs in which favorable adhesion between diffusion barrier layer and the IMD layer cannot be achieved. For example, copper lines peel off at the interface, worsening the yield of the semiconductor device.\nTherefore, there is a need for an improved diffusion barrier layer in the copper interconnect, and a method of forming thereof."}
{"text": "The invention relates to a battery terminal connecting cable with a strand consisting of numerous fine wires and a further section having a hole for a screw, whereby on the further section the numerous fine wires are welded together, a method for the manufacture of a battery terminal connecting cable and a device for the manufacture of such a cable.\nBattery terminal connecting cables are used primarily to interconnect battery cells conductively. In this case one also talks of cell connectors. For the conduction of relatively high currents at low electrical resistance the strand usually consists of copper wires twisted into a strand having a cross-section of approximately 50 mm2. Depending on the field of application, however, both thicker and thinner strands are used.\nIn order to connect such a strand consisting of numerous fine copper wires to a battery terminal, a copper pipe section is first inverted over the strand and this pipe section is then pressed to an approximate plate shape. In this plate there is a hole which first passes through the upper side of the original pipe section, then through the compressed cable and finally through the lower side of the original pipe section. Finally in this hole there is inserted a screw which interacts with a thread in the battery terminal so that when the screw is tightened, the strand held together by means of the copper sleeving is pressed onto the battery terminal.\nIn practice it has been found that even a relatively strongly tightened screw works loose as a result of vibrations, as occur especially in batteries located in vehicles. This has the result that the end of the cable is no longer securely connected to the battery terminal. Contact is therefore made over smaller areas and is thus severely impaired. In cases of more severe loosening this leads to undesirable heating, even to the extent of sparking.\nIn order to prevent screws working loose, in known battery terminals a plastic section, such as preferably a plastic bead, is provided in the thread region between the screw and the battery terminal. This bead becomes deformed as the screw is turned and acts to prevent loosening of the screw.\nThe insertion of a bead or other plastic section has the result that expensive, special screws are required. Since battery terminal connectors are a mass-produced product manufactured in large quantities, any increase in the cost of the product leads to economic disadvantages.\nA battery terminal connecting cable of appropriate type is known from EP-A-O 707 321. In practice the end of the cable strand produced by spot welding does not give the required strength values for the cable end so that the cables described have not proved suitable as battery terminal cables.\nU.S. Pat. No. 4,325,760 discloses a cable whose end was soldered. However, soldering is an expensive procedure which leads to increased costs in the area of mass production.\nThe invention is based on the problem of preparing a connection between battery terminal connecting cable and battery terminal at a favourable cost, which ensures secure contact between the battery terminal connecting cable and the battery terminal even under severe vibration or thermal expansion.\nThe problem is solved using a battery terminal connecting cable of appropriate type where the strand is at least partly insulated and a further section is secured at the end of the strand.\nThe invention is based on the knowledge that the problem of screw loosening in the battery terminal is not primarily attributable to the screw being held inadequately in the battery terminal. The reason for loosening screws lies in the fact that at times fairly high temperatures occur at the contact piece of the battery terminal connecting cable as a result of the flowing currents. At varying temperatures the fine wires inside the copper sleeving are severely compressed from time to time and then remain in this severely compressed form. As a result after many temperature fluctuations the originally tightened screw is subject to less force so that the screw sits more loosely in the thread. Loosening of the screw can then only be impeded by the plastic section described above. The plastic section certainly impedes loosening of the screw but the fact that the screw at times exerts less pressure on the contact piece is not impeded by the plastic section. However, a constant pressure on the contact piece is required to ensure a continuous equally good current flow.\nPractical tests have shown that when a cable end consisting of wires welded together, is screwed on, the compression of the connecting piece described above does not occur and consequently normal screws without protection against loosening can be used. Since the screw now presses the contact piece onto the battery terminal at a constant pressure, optimal electrical conditions are created and thus the heating in the contact region is reduced.\nAccording to the invention a further section is secured at the end the strand. This further section can, for example, serve as a contact piece and can be connected in various ways to the end of the strand welded to form a solid material.\nA particularly preferred embodiment of a battery terminal connecting cable is obtained if the further section is welded onto the end of the strand. In this way a specially favourably shaped contact piece can be used. This in particular allows the weight of the battery terminal connecting cable to be reduced.\nIt is also advantageous if the further section, preferably constructed as a contact piece, is made of copper.\nIn all known battery terminal connecting cables the longitudinal axis of the further section extends in the direction of the axis of the battery terminal connecting cable. Since the battery terminal connecting cable serves as a connecting element between two terminals, it seems to be logical to use a further section extending in the direction of the axis of the battery terminal connecting cable. However, tests have shown that in many cases it is of great advantage if this further section is configured such that its longitudinal extension runs at an angle to the longitudinal axis of the cable. Examples of this type of configuration are shown in FIGS. 6-9.\nDepending on the configuration of the battery terminals, in practice the connecting cables are frequently very severely bent. In many cases, the angular configuration of the further section to the axis of the battery terminal connecting cable reduces the necessary bending of the cable and also allows shorter cables to be used. The configuration described thus has the result that the cable can be protected and material usage can be reduced.\nThe angular configuration described is also of great importance for battery terminal connecting cables regardless of the afore-mentioned characteristics and can thus also be used to advantage for clamp connections between the cable strand and a further section.\nThe problem according to the invention is also solved using a method for the manufacture of a battery terminal connecting cable in which a strand consisting of numerous fine wires is welded by means of ultrasound and the strand is welded to a contact piece or with a contact piece.\nThe use of ultrasound for making electrical connections is already known from EP-A-0 671 790. However, the welding of the strand to a contact piece or with a contact piece cannot be inferred from this document.\nAll types of welding or soldering of an end of a strand developed so far have failed in that the high heat input during the welding process is transferred from the copper cables normally used to the insulating layer. On the one hand, this leads to destruction of the insulation and on the other, the input heat is rapidly dissipated. Only welding by means of ultrasound produced such good results that the insulating layer can remain on the cable strand during welding. This surprising result reduces the energy input and allows known manufacturing methods for moulding plastic onto the connecting regions to be used even for the new battery terminal connecting cable.\nFor manufacture of the battery terminal connecting cable according to the invention a device is proposed which has a sonotrode, a strand feed device, a contact piece feed device and a pressure cylinder arranged so that the strand and the contact piece can be pressed onto one another by means of this device.\nThe device allows the battery terminal connecting cable to be manufactured automatically and rapidly.\nSince during pressing of the wire strand the fine wires are pressed to the side, it is proposed that on the device according to the invention there should preferably be provided movable jaws which act on the strand at right angles to the axis of the pressure cylinder. These jaws together with the pressure cylinder and the sonotrode define a cross-section in which the wires should be welded together to form the most compact copper piece. At the end of the welding process the jaws are moved away from the strand and the pressure cylinder is retracted so that the strand connected to the contact piece can be removed from the device."}
{"text": "As known in the art, the CD (compact disc) system currently employed as a main-current acoustic apparatus is of such a type that pits are formed in a 12 cm-diameter disc in a manner to correspond to digital PCM (pulse code modulation) data and reproduction is performed by, while rotating the disc in a drive rotation in a CLV (constant linear velocity) system, linearly tracking the data from an inner periphery side to an outer periphery side on the disc by a semiconductor laser and an optical pickup built in a photoelectric conversion element.\nAudio playback musical tone data and sub-code data P, Q, R to W defined for control/display are recorded in the disc. Of the aforementioned data, the sub-code Q is also called as address data and represents, in a program area containing disc's musical tone data, a program number (TNO) of recorded musical tone data, a phrase number (INDEX), a lapse time (TIME) for each program, a total time (ATIME) lapsing from a start position of the program area, etc.\nFurther, the aforementioned sub-code data Q represents, in a lead-in area located on the inner periphery side of the program area, a start address of each program as TOC (table of contents) data. That is, the sub-code data Q is recorded for implementing a search operation, accurately at high speeds, for selective reproduction of playback data from vast information recorded in the disc to enable a stereo reproduction to be performed for about one hour and grasping the state of a reproduction on the disc.\nThe CD system has initially been developed so as to record and reproduce musical data. In recent years, with attention paid to such a vast recording capacity, a CD-ROM (read only memory) system is determined to be used as a read-only data recording medium for the disc through the use of a musical data recording area of the disc as a digital information recording area. The CD-ROM system is of such a type as to record and reproduce digital information on the disc without varying a recording/reproduction format for musical data reproduction in the CD system and to do so by adding a new format thereto.\nIn the existing digital data processing system using a CD-ROM system, however, when there occurs a shift in frequency of a clock serving as a reference of a data processing speed relative to a host processing section side and disc reproduction section side, the buffer memory at the host processing section side is placed in an over- or an underflowed state, failing to make a normal data reproduction and thus presenting a problem.\nThe present invention is achieved, taking the above situations into consideration and the object of the present invention is to provide a much better digital data processing apparatus which, even when there occurs a shift in frequency of a clock serving as a reference of a data processing speed relative to a host processing section side and disc reproduction section side, can prevent a buffer memory on the host processing section side from being placed in an over- or an underflowed state and make a normal data reproduction."}
{"text": "1. Field of the Invention\nThis invention relates generally to tunnel boring machines, and more particularly to a machine that utilizes a casing centered within a casing structure, and still more particularly to one which utilizes a spoils ejector adjacent the cutting head at the forward end of the auger that is turned by the boring machine.\n2. Description of the Related Art\nThis invention is directed to the providing of an auger boring machine of the type used in the installation of in-ground pipe. For decades, prior art boring machines have been positioned on tracks in a large hole at the initiation site of a tunnel. The machine has a power source that pushes a casing into a tunnel dug by a rotating drill head. This digging is accomplished by the power source turning an auger that extends within and is parallel to the casing. At the terminal end of the auger is a drilling head that revolves as the auger spins, the drilling head cuts soil and/or rocks loose from the ground. The cut soil and/or rock are referred to as spoils and are evacuated from the tunnel by the auger blade.\nPrior art tunnel boring machines bore tunnels 4-96″ in diameter. The machines cost thousands of dollars to set up to begin the digging of a tunnel. When adjustments need to be made, it can take a day. Oftentimes, the auger used in these machines has to be removed from within the casing if a problem is encountered at the cutting head. This takes time and incurs expense.\nIt is thus apparent that the need exists for an improved tunnel boring machine which in addition to being easy to use, permits tunnels to be bored faster and more economically than in the past."}
{"text": "1. Field\nThe present disclosure relates to a composition for preventing, improving, or treating osteoporosis (Type I osteoporosis) in postmenopausal women, the composition including, as an active ingredient, scopolin, a derivative thereof, or a pharmaceutically acceptable salt thereof.\n2. Description of the Related Art\nEstrogen which is a representative hormone that symbolizes women is an important female hormone that regulates a lifespan of women through menstruation, pregnancy, and menopause. Estrogen is a well-known female hormone that is mainly secreted from follicles and corpus luteum that are in female ovaries and from the placenta, and is a generic term for estrone (E1), estradiol (E2), and estriol (E3). Estrogen affects a wide range of tissues, and is especially needed to maintain flexibility and steady state of the uterus, urinary tract, breast, skin, bones, and blood vessels. In particular, estrogen acts on osteoblast involved in production of bones to maintain bone formation, and also acts on osteoclast involved in removal of bones to inhibit bone resorption, thereby maintaining homeostasis of bones (see FIG. 1). Thus, postmenopausal osteoporosis which is caused by reduced estrogen levels is one of the most serious diseases in postmenopausal women (see FIG. 2).\nOsteoporosis is a condition where bone mass per unit volume is abnormally reduced compared to a normal level according to sex, age, and race of a normal person. Due to brittle bones that break easily by minor impacts such as bending the waist or sitting down, osteoporosis is characterized by fractures in the coxa, carpal, vertebrae, or the like. According to pathogenesis of osteoporosis, osteoporosis is divided into primary and secondary forms. Primary osteoporosis includes postmenopausal osteoporosis (Type I osteoporosis) and senile osteoporosis (Type II osteoporosis), and secondary osteoporosis is caused by drugs or the like (see FIG. 3).\nRegarding postmenopausal osteoporosis (Type I osteoporosis) in women, when a deficiency of estrogen hormones appears after menopause, postmenopausal osteoporosis occurs as components of bones are absorbed into body tissues and calcium absorption through intestines is reduced (see FIG. 4). In addition, the reduction in estrogen hormones inhibits differentiation and proliferation of osteoblasts, resulting in inhibition of bone formation, and activation of osteoclasts leads to occur bone the resorption more frequently than the bone formation, resulting in an increase of bone loss and a decrease of bone mass (see FIG. 5). The above phenomenon may progress rapidly after menopause depending on a person.\nSenile osteoporosis (Type II osteoporosis) is caused by bone loss with increasing age in both men and women. In particular, senile osteoporosis is characterized by a decrease in intestinal calcium absorption due to a reduction in active vitamin D in the body and by a decrease in the number of osteoblasts that newly produce osteocytes. Here, senile osteoporosis progresses relatively slowly herein.\nSecondary osteoporosis is caused by various diseases or medications that affect functions related to bone cell production and maintenance of the human body. For example, such various diseases may include hyperthyroidism, hyperparathyroidism, Cushing's syndrome, rheumatoid arthritis, hyperprolactinemia, and the like, and in addition, steroid hormone preparations or thyroid hormone preparations may also cause secondary osteoporosis.\nWhen the climacteric begins, female hormones secreted from the ovaries are reduced due to a decrease or an imbalance in ovarian functions, and accordingly a variety of symptoms may appear. Even after menopause, postmenopausal osteoporosis may occur due to a rapid decrease in bone density and a decrease of bone mass. Since postmenopausal osteoporosis may cause backache or other bone-related diseases, and bone fractures, postmenopausal osteoporosis is considered as a disease that significantly affects the quality of life. In particular, women undergoing premature menopause or having an ovariectomy before age of 50 may be more susceptible to postmenopausal osteoporosis.\nAs therapeutic agents for postmenopausal osteoporosis in women, incrementally modified composite drugs, such as Bonviva Plus (composition: bisphosphonate-based drugs) and Aclasta (composition: zoledronic acid injection 5 mg), may be used. However, due to side effects of the agents above, the demand for therapeutic agents derived from natural products or health functional foods is increasing. Calcium which is good for bone health, vitamin D, and health functional foods including isoflavones are widely used, but effects thereof are limited. Therefore, there is a need for the development of more effective pharmaceutical preparations or health functional foods for prevention, improvement, and treatment of postmenopausal osteoporosis in women."}
{"text": "1. Field of the Invention\nThe present invention relates to data processing systems, and more particularly to a data processing system, method, and computer program product for maintaining a relational database in a virtualization system of usage information of multiple, different physical devices in a system that implements virtual storage devices.\n2. Background of the Invention\nIn some known data processing systems, a server computer system is coupled directly to its storage devices, such as tape drives and/or disk drives. To perform a backup operation to a tape drive, for example, the server computer system commands the tape library to mount a specific physical tape cartridge into a specific physical tape drive. The application running on the server can then directly access the physical tape drive. Any warning or errors encountered can be managed by the server through the application. It is possible for this error information to be used by the application to perform preventative maintenance on the physical storage devices. The problem with this environment is that the storage administration resides on the server. The server has access to only its own storage devices. These storage devices cannot be shared with other servers.\nStorage virtualization has been used to overcome the limitations that arise when a server is coupled directly to its storage devices. In a system that implements storage virtualization, the server computer systems are isolated from the physical storage devices by a virtual representation of the storage system. The storage system is a pool of physical storage devices, such as disk drives and tape drives. This approach permits the sharing of the physical storage devices among the server computer systems. Each server has access to one or more virtual storage devices. These virtual storage devices are logical devices maintained by the virtual interface. The virtual interface maps the virtual storage devices to physical devices. When an application running on one of the servers wants to access data on ore of its virtual interfaces, the server identifies the virtual interface. For example, when a server wants to perform a backup, the application commands the tape library to mount a physical tape cartridge into a virtual tape drive. The virtual interface then intercepts the command, modifies the request to translate this virtual tape drive to a physical tape drive, and creates a mapping from this physical device to the virtual interface.\nStorage virtualization hides warnings and errors from the server, however. The server knows about only the virtual interface. Usage information about the physical device is not available to the server. Although, the virtual interface logs some limited information about what occurs with the physical storage devices, the log is a continuous event log that combines very minimal information about events that occurred in the physical devices. The log is a single, sequential string of events occurring in various physical devices kept in the log in the sequential order in which the events occurred. Events are placed in the log in the order in which they occur without regard to where the event occurred. Thus, an event occurring in a first storage device may be logged followed by an event occurring in a second storage device followed by an event occurring back in the first storage device. It does not maintain a separate log for each device. Further, information about which physical device was used and the history of usage is not recorded. This log is not available to the servers. They may not access it.\nThe StorageNet 6000, available from Storage Technology Corporation, is an example of a system that implements a virtual interface. The SN6000 presents a virtualization system for tape drives. The system presents virtual tape drives to servers, and dynamically maps these virtual interfaces to the physical tape drives as needed. When errors occur in the system, it is very difficult for the server to determine where the error occurred since it does not have information about which physical device was being used. When an error occurs, a host attributes the error to the virtual interface that the host is coupled to, and not the physical device.\nTherefore, a need exists for a method, system, and product for maintaining a relational database within the virtual interface of information about usage of physical storage devices in a system that implements virtual storage devices."}
{"text": "One type of roller conveyor is powered by DC motors located in roller cylinders which are mounted with other rollers between the conveyor side plates. Electronic motor control modules are mounted on one of the side plates. Packages traveling on the conveyor rub against the side plates and rollers generating static electricity which can arc between the side plate and control module damaging the control module components."}
{"text": "Traditional editing for storytelling on film (and later video) created and perfected the style of inter-cutting shots of varying focal lengths and at varying angles to communicate the positional relation between characters in establishing shots and to provide increased intensity of performance by using close-ups. The ability to cut away from one perspective to another also allowed the use of different takes within the same production. These techniques have become so commonplace that they are inculcated and accepted by everyone with any media exposure whatsoever.\nThese techniques have also been used for many years in live productions with multiple cameras from the beginning of Television as Live Broadcast TV to today's news, awards, sports and talk shows.\nSince the mid 1960's as the “Studio Tours” flourished in the Los Angeles area, “simulated” shows were created where participant could take the role of actors on sets with contemporary equipment and create short video scenes from famous shows. This type of “live” production has become a Theme Park staple over the years and many versions of custom media inserting guests into the previously produced segment have been created.\nOn a separate path, beginning more than 20 years ago, as MTV exploded onto the scene, the evolution of audio Karaoke to music video setups were taking the first steps of simplifying the production studio apparatus to bring some version of these expensive production techniques to a new market as a simplified production studio. These were all single camera set-ups and were not attempting to accomplish anything more than a superimposition of the guest over an interesting background while the music played. There was no intent of interacting with the background.\nThe next step in the evolution was to insert participants into a background entirely pre-produced specifically for the purpose of inserting a participant member of the public who could interact by answering questions or reading prompts to give them the appearance of appearing on a television show. In one example, this was showcased a decade ago at EPCOT center in Florida where guests could do a brief interview with Jay Leno. This setup was reprised at the Olympic Village in Atlanta in 1996 and has been resident at the NBC Experience in Rockefeller Plaza since 1999. There are other options in the later location and another example in “Studio 39” in San Francisco. However, in all of these installations, a single camera is used to place the participant into a back ground scene produced specifically and only for that purpose."}
{"text": "In current aircraft technology there is the hydraulic system for operating inter alia the flight control system. This hydraulic system is divided into three different circuits (green circuit, yellow circuit, blue circuit). Under high loads in these circuits the maximum allowable temperature of the hydraulic fluid (e.g. Skydrol) of 90° C. can be reached, which causes high stress in the individual components (pumps, conduits, and seals). This reduces the reliability of these components.\nA fuel cell is an electrochemical reactor, which produces electric and thermal energy from the reactive gases hydrogen and oxygen. The electric energy is conducted from the cell e.g. as direct current and conditioned accordingly for further use. The thermal energy, which is produced, has to be removed through a cooling circuit, so that the fuel cell does not overheat during operation, and does not get damaged."}
{"text": "Fast dissolving pharmaceutical dosage forms which are designed to release an active ingredient in the oral cavity are well known and can be used to deliver a wide range of drugs (Critical Reviews in Therapeutic Drug Carrier Systems, 21(6):433-475 (2004); Seager H. (1998), J. Phar. Pharmacol 50:375-382; Bandari et al. (January 2008), Asian Journal of Pharmaceutics 2-11).\nIn a fast dissolving dosage form, a drug may physically be trapped in a matrix composed of e.g. mannitol and fish gelatin (EP 1501534; EP1165053), modified starch (U.S. Pat. No. 6,509,040), pullulan in combination with an amino acid (EP1803446), or maltodextrin in combination with sorbitol (US2004/0228919). The solution, suspension or dispersion of the drug and the carrier material may be filled into blister cavities, frozen and thereafter lyophilized. As any product of technology at all times, also the latter can still be improved upon.\nThe subject invention provides an improved fast-dispersing dosage form."}
{"text": "1. Field of the Invention\nThis invention relates to the production of hydrogen-rich gas by the partial oxidation of a water saturated gaseous hydrocarbonaceous fuel or a mixture of water saturated gaseous hydrocarbonaceous fuel and a liquid hydrocarbonaceous fuel.\nThis invention also relates to the use of at least a portion of the hydrogen-rich gas or synthesis gas produced during a partial oxidation gasification process as the feedstock in a process for the direct reduction of iron, also referred to as the \"DRI\" process.\n2. Description of the Prior Art\nThe production of hydrogen-rich gas by the partial oxidation of a liquid hydrocarbonaceous feed or a pumpable slurry of solid carbonaceous feed to produce a raw feed gas which is cleaned and purified to produce a product gas comprising about 20 to 60 mole % hydrogen is disclosed in coassigned U.S. Pat. No. 4,338,292. In coassigned U.S. Pat. No. 4,021,366, the process gas stream from the partial oxidation of a liquid hydrocarbon material or a slurry of solid carbonaceous feed is cooled, scrubbed, and water-gas shifted in a bed of low temperature catalyst followed by a bed of high temperature catalyst. In coassigned U.S. Pat. No. 5,152,976, refinery offgas is partially oxidized.\nPartial oxidation gasification processes are most effective when operated at relatively high pressures, generally greater than 20 atmospheres. As used herein, a partial oxidation reactor can also be referred to as a \"gasification reactor,\" or simply a \"gasifier\" and these terms are often used equivalently and interchangeably.\nThe feedstock for a partial oxidation reaction is usually a hydrocarbonaceous material, that is, one or more materials, generally organic, which provide a source of hydrogen and carbon for the gasification reaction. The hydrocarbonaceous material can be in a gaseous, liquid or solid state, or in a combination as desired, for example, a solid-liquid composition in a fluidized state. Petroleum based feedstocks include petroleum coke, coal residual oils and byproducts from heavy crude oils. The coal or coke can be in a finely divided state. Waste plastic materials can also be used as the feedstock.\nMany uses for the syngas produced from the partial oxidation reaction are at relatively lower pressures. Therefore, expansion of the high pressure syngas across a power recovery machine to obtain a syngas at reduced pressure is often employed. This type of expansion is often used as a means for generating electricity. The power generation step is not 100% efficient and some energy is lost in converting the energy to electricity. The electricity generated in such an expansion process requires step up transformers, additional electrical switchgear, and a use for the electricity.\nThe DRI process can use syngas as a feedstock if it is fed to the reaction chamber at a relatively low pressure, typically less than about 5 atmospheres for a moving bed reactor, and less than about 15 atmospheres for a fluidized bed reactor. The DRI offgas is cooled, compressed, sent to a carbon dioxide removal step, and then mixed with fresh syngas feed and recycled to the DRI process. The recycle compressor is a large energy user and often uses an electrically driven motor. This electrically driven motor is not 100% efficient and some energy is lost in converting electricity to shaft power. Use of an electrical motor requires step down transformers, additional electrical switchgear, and a source of electricity. Alternately, a steam driver can be used with similar energy losses and supporting equipment.\nThe operation of gasifiers at relatively high pressures in power generation systems, for example, the integrated gasification combined cycle (IGCC) system is disclosed in U.S. Pat. Nos. 5,117,623 and 5,345,756, wherein these systems are coupled with expanders, gas turbines and steam turbines for power generation. U.S. Pat. Nos. 5,531,424 and 5,370,727 disclose processes for the direct reduction of iron."}
{"text": "1. Field of the Invention\nThis invention relates to the field of sonar, and more particularly to the field of demultiplexer circuits used in demultiplexing the multiplexed signal information obtained from DIFAR sonobouys.\nDIFAR sonobouys are dropped in numbers by intelligence gathering aircraft into ocean areas that are to be subjected to short term sonar monitoring. As shown in FIGS. 1a and 1b respectively, the DIFAR sonobouys typically have a first dipole directional sonar hydrophones that respond to sound along a sonobouy (Y) axis, a second dipole directional hydrophones that respond to sound along a sonobouy (X) axis, and a non-directional or omni hydrophone that responds to sound with uniform sensitivity on all level radials from the instrument.\nThe DIFAR sonobouy is typically characterized to float at a predetermined depth for a predetermined duration. The (X) and (Y) axes are orthogonal and are kept float level by operation of the hydrostatic design of the sonobouy. The sonobouy is equipped with a fluxgate compass. No means is provided to orientate the sonobouy once it is in the water.\nSounds received by the hydrophones in the sonobouy are multiplexed in accordance with a predetermined DIFAR format for transmission to an aircraft sent to record the received data.\nThe DIFAR format requires that the following conventions be adopted. Using a top view of the sonobouy, angle Phi is measured clockwise from magnetic north to the +Y axis. Angle Theta is measured clockwise from the +Y axis to the target location. When dropped into the ocean, the angle phi can be any angle from 0 to 360 degrees. The sonobouy orientation can continually change due to wave action. The relative amplitude and polarity of target signals from the Y axis hydrophones varies as the cosine of angle theta. The relative amplitude and polarity of target signals from the X axis hydrophones varies as the Sine of angle theta. The omni hydrophone signal is positive in all quadrants of angle theta.\nThe sonobouy circuitry provides a 15.0 kilohertz phase reference signal. The 15.0 kilohertz phase reference signal is phase shifted or delayed in time by an interval proportional to the magnetic heading angle phi to generate a signal referred to as the phase pilot signal.\nThe Y axis hydrophone information is referred to as cosine channel information. The X axis channel information is referred to as sine channel information. The cosine channel information is multiplied by a reference carrier signal having a frequency of 15.0 kilohertz. The reference carrier signal leads the 15.0 Kilohertz phase reference by 90 degrees. All delays are measured with respect to the phase reference signal. The sine channel information is multiplied by a quadrature reference carrier signal having a frequency of 15.0 kilohertz and is phased shifted to lead the reference carrier of the cosine channel by 90 degrees.\nThe sonobouy also provides a constant amplitude 7.5 Kilohertz frequency pilot signal. The frequency pilot signal is by definition exactly one half the frequency of the phase reference signal.\nThe omni signal is not multiplied by a carrier. The bandpass of the omni signal information is typically under 2.5 kilohertz. An RF carrier is modulated with the sum of the frequency pilot, the phase pilot, the cosine channel, sine channel and the omni signal information and transmitted by the DIFAR sonobouy to the sampling aircraft to be recorded on board the aircraft or for retransmittal to a receiving station as a single channel signal for demultiplexing by the invention DIFAR demultiplexing circuit.\n2. Prior Art\nPrevious DIFAR demultiplexer circuits used two or more phase locked servos. The invention DIFAR demultiplexer circuit uses one phase locked servo loop with a digital correction in 45 degree steps and analog adjustment of the quadrature phase signal to the phase pilot."}
{"text": "1. Technical Field\nThe present invention relates to a method of upgrading low-value linear alpha olefins (LAOS) of C18 or more. More particularly, the present invention relates to a method of manufacturing high-quality lube base oil by converting a low-value LAO having a C18 or more hydrocarbon chain into a linear internal-olefin (LIO) having a controlled double bond position distribution, which then undergoes oligomerization.\n2. Description of the Related Art\nEthylene is a raw material used in the chemical industry to the extent that the production and consumption thereof are considered as indicators of the chemical industry scale of a country. Typically, ethylene is employed as a monomer for preparing a polymer such as polyethylene. In some cases, LAOs having a C4-C40 carbon length (or chain) are prepared by adjusting the degree of polymerization of ethylene, and thus utilized to manufacture a variety of chemicals. In this regard, as used herein, the term “full range LAO” refers to a C4-C40 LAO mixture resulting from ethylene polymerization or oligomerization.\nThe reason why LAO production techniques are important is that LAO is a chemical without sulfur and nitrogen derived from crude oil. Generally, since crude oil contains impurities such as sulfur and nitrogen components in amounts of ones of wt %, it is difficult to directly prepare chemicals composed of pure hydrocarbons without such impurities. However, ethylene resulting from catalytic cracking of crude oil is converted into LAO, which is then converted again into a desired chemical, thereby obtaining chemicals composed of pure hydrocarbons without impurities.\nEthylene polymerization is mainly carried out using a batch reactor under air-sensitive conditions in the presence of a metallocene catalyst. Since a metallocene catalyst has a very strong single acid site, ethylene is selectively polymerized at the single acid site of the catalyst, and is thus linearly grown at the acid site of the catalyst. Ultimately, the olefin produced by the LAO preparation reaction possesses a linear alpha-olefin structure having an even number of carbon atoms because of the use of an ethylene monomer.\nWhen operating conditions for the LAO preparation reaction are controlled, the LAO composition distribution may be altered. For example, when the reaction temperature or the polymerization time is increased, the proportion of LAO having a relatively high molecular weight may become high. In contrast, when the reaction temperature or the polymerization time is decreased, the proportion of LAO having a relatively low molecular weight may increase. Taking into consideration the industry scale for LAO-derived chemicals by LAO manufacturers, the LAO composition is changed in such a manner that the reaction conditions are appropriately controlled, thereby ensuring optimal profitability for individual projects.\nAlpha-olefins produced by the LAO preparation processes exhibit properties varying depending on the number of carbon atoms thereof, and the properties of the resulting chemicals may become different depending on the type of LAO feed. For example, C4 LAOs obtained by polymerizing two ethylene monomers are present in a gas phase, and polymers resulting from the polymerization thereof have an excess of intramolecular branches, making it difficult to apply to a specific product, for example, lube base oil. Also, even when C6 LAOs obtained by polymerizing three ethylene monomers are converted into polymers, such polymers are unlikely to be applied to lube base oil and so on due to many intramolecular branches thereof. As such, copolymerization of ethylene with a C6 LAO may result in an ethylene-1-hexene copolymer having properties different from those of conventional polyethylene. Also, C8 LAOs obtained by polymerizing four ethylene monomers may be applied to Group IV lube base oil via polymerization, and may be used to produce a copolymer with ethylene, as in the C6 LAOs. Also, C10-C12 LAOs are polymerized and thus mainly used as Group IV base oil, and C14-C16 LAOs are reacted with amine or succinic acid, so that the resulting products may be applied to various functional chemicals, or mixed and thus used for cheap drilling oil.\nThe aforementioned LAOs have certain end-uses, and the markets adapted for the demand thereof are established. Furthermore, the related techniques have been already developed to some extent. The most expensive LAO is a C8-C12 LAO for preparing high-quality lube base oil (e.g. Group IV lube base oil), and C6 and C8 LAOs may be mainly utilized because the demand thereof is high for producing copolymers. Also, although C14-C16 LAOs are relatively less used due to limited end-uses thereof compared to the C6-C12 LAOS, C14-C16 LAOs may be used to prepare specialty chemicals and may thus be formed into products by increasing purity thereof. However, C18 or more LAOs are problematic because the end-uses thereof are limited, and the purity thereof has to be increased to obtain desired products. In this regard, C18 or more LAOs have been employed as an additive for lube oil or in wax form, but are not general-purpose chemical products, and thus the demand thereof is undesirably low. Also, C20 or more LAOs are limited in serving as feeds because they exist as a solid at room temperature. Like this, the market prices of C18 or more LAOs are maintained very low compared to the other LAOS, on account of low demands, tendencies to exist as a solid, and restrictions where such feeds may be prepared into products only when the purity thereof is increased.\nOwing to the limitations of end-uses as above, the yield of high-value C8-C12 LAOs is intended to be selectively increased while decreasing the yield of C18 or more LAOs, by LAO manufacturers. However, the LAO products have a composition of the normal distribution due to characteristics of the polymerization catalyst, making it impossible to attain operating conditions for preparing only C8-C12 LAOs and for drastically lowering the yield of C18 or more LAOS.\nHence, LAO preparation under reaction conditions able to lower the yield of C18 or more LAOs as much possible is chosen by LAO manufacturers, and is advantageous because only LAOS, the demand of which is high and the value of which is also high, may be selectively prepared while decreasing the yield of C18 or more LAOs regarded as byproducts. However, since the composition of products resulting from ethylene polymerization has the normal distribution as mentioned above, the yield of C8-C12 LAOS, corresponding to feeds for end-uses (especially Group IV lube base oil) which have the highest price and a future market of which is expected to grow, is relatively low, which is undesirable.\nAlternatively, the reaction conditions may be set so as to maximize the yield of the most expensive C8-C12 LAOs without taking into account the yield of C18 or more LAOS. As such, the yield of C18 or more LAOs is 10 wt % or more, but the yield of relatively high-value C8-C12 LAOs may be increased. Therefore, if C18 or more LAOs obtained by full range LAO techniques can be simply converted into general-purpose high-value products, the effect of such technology is expected to be very significant. Accordingly, methods of converting C18 or more LAOs into high-value chemicals are required.\nMeanwhile, techniques for preparing chemicals (e.g. drilling oil) having improved properties in such a manner that a double bond present at the alpha position of the LAO is shifted to the inside of the carbon chain to form a linear internal-olefin (LIO) (olefin shift reaction) are also known (e.g. U.S. Pat. No. 6,281,404). However, the above techniques are not regarded as the application of low-value C18 or more LAOs because a wide range of C4-C50 LAOs are utilized, and also because the double bond shift reaction for 1-pentene is mainly disclosed. In addition, there are known techniques for manufacturing lube base oil (e.g. EP 0 136 377 B1) by oligomerizing internal olefins (especially, internal olefin produced from dehydrogenation of paraffin so that a double bond is randomly distributed in the molecule) in the presence of a boron trifluoride catalyst and a promoter, and for manufacturing lube base oil (Japanese Patent Application Publication NO. 1993-306241) by oligomerizing a C10-C24 LAO/LIO mixture. However, these techniques are limited in improving properties of the lube base oil because the position of the double bond in the olefin is not controlled. In particular, such techniques are unlikely to become technology for increasing the value of low-value C18 or more LAOS produced inevitably as byproducts in full range LAO preparation."}
{"text": "The present invention relates to a high-frequency amplifier circuit that outputs differential signals, a semiconductor device, and a magnetic recording and reproducing device.\nHigh-frequency amplifier circuits used in radio communication and the like widely employ a multistage configuration using a plurality of amplifiers so as to attain desired power for a load. Furthermore, differential signals are used for various purposes, such as improving noise immunity and reducing power consumption, and in cases where it is difficult to place a ground electrode. However, in the case where a high-frequency amplifier circuit having a multistage configuration outputs differential signals, for example, when a common-mode component is generated in the circuit, abnormal oscillation may occur in the amplifier via a feedback path, such as a power line or a signal line, which has been a problem. Note that in order to output differential signals at the last stage in a high-frequency amplifier circuit, one differential amplifier including two amplifiers having the same properties or a dual configuration including two amplifiers having the same properties placed in parallel may be used."}
{"text": "Cigarette lighters have usually been independently carried in pockets and as a result were often inadvertently left on a table by the owner or momentarily loaned to another who thoughtlessly carried it away. No satisfactory way of securement has been conceived for attaching the cigarette lighter to the wearing apparel for avoiding such occurrences and which also provide an attractive apearance and enabled the wearer to use the lighter without the necessity of separating the lighter from its securement of the wearing apparel."}
{"text": "The present invention generally relates to stringed musical instruments and more particularly to a stringed instrument, such as an acoustic guitar, comprising a sound chamber, or body, in which sound waves generated by the plucked strings are amplified by the vibrations of the materials forming the sound chamber and emitted from the sound chamber. For musical instruments, such as flat-topped acoustic guitars, the sound chamber has a front (also referred to as the soundboard), back, and sides. For some of these instruments, the strings are attached on one end to a head at the end of a neck extending from the sound chamber, and attached at the other end to a bridge which is attached to the soundboard. The flat-topped acoustic guitar is so identified because its soundboard is generally flat, as opposed to an arched top guitar which, as suggested by the name, has an arching top with three dimensional characteristics. The term flat-topped acoustic guitar includes both steel string acoustic guitars as well as nylon string classical guitars.\nUnder the traditional design of a flat-topped guitar, the soundboard has a single opening for the emission of sound waves, which opening is disposed beneath the strings of the instrument, the opening being located centrally with respect to each side of the instrument.\nWhen the strings vibrate, the vibrations travel through the bridge to the soundboard, such that the entire soundboard vibrates. The rest of the sound chamber amplifies the vibrations of the soundboard, with the sound waves primarily emitted from the soundhole.\nIn the case of acoustic guitars, the sound chamber generally comprises a top, known as the soundboard, a back, and sides. The lower bout is the large rounded bottom of a traditional guitar and the upper bout is the smaller, rounded and convex shape at the top. Under traditional design, the shape of the sound chamber of acoustic guitars is in the shape of a number “8”, with the upper half, i.e., the upper bout, being smaller than the bottom half, i.e., the lower bout. The upper bout and lower bout are separated by the “waist” of the guitar, which is the concave transition between the upper and lower bouts. For traditional designs, the sound hole is located almost entirely in the upper bout of the instrument.\nFor the typical right handed guitar player, the upper bout of the guitar is adjacent to the player's left arm, and the lower bout is adjacent to the player's right arm. The left hand is utilized for fingering notes on a fingerboard (also referred to as the “fretboard”), where the fingerboard is disposed on the neck. The right hand is utilized for picking or strumming the strings. For the remainder of this description, it will be assumed that the guitar is a “right handed” guitar, i.e., built to be played by a right-handed person. However, correlating the description for a left-handed guitar only requires the assumption that the right hand is utilized for fingering the notes and the left hand is utilized for picking or strumming the strings.\nThe sides of the guitar may, for purposes of description, be identified with respect to the strings. The treble strings of the instrument are on the side of the instrument facing downward as it is played in the conventional manner, while the bass strings are on the side of the instrument generally facing upward as the instrument is played. Using the strings as a point of reference, the sides of a guitar may be referred to as the treble side and the bass side. With respect to the string orientation described above, the side of the guitar facing downward while being played is considered the treble side of the instrument and the side of the guitar facing upward is considered as the bass side of the instrument. The upper bout therefore may be further described as having a treble side upper bout and a bass side upper bout. Likewise, the lower bout may be further described as having a bass side lower bout and a treble side lower bout. The bass and treble sides of the lower bout meet at the bottom center of the guitar, where the meeting sides of the lower bout are glued to an an internal structure known as the “tailblock.”\nUnder the traditional design for guitars, the exterior of the sound chamber is symmetric, where the treble side and bass side are mirror images of one another. However, over the years, instrument makers have modified the traditional design. One of the most common of these modifications, which results in an asymmetrical sound chamber, has been to fashion a “cut-away” into the treble side of the upper bout and upper portion of the soundboard adjacent to the neck on the treble side to allow the player greater access to the portions of the fingerboard adjacent to the body of the guitar. Other modifications have also been made, such as placing the sound hole in a different position than directly under the strings. As another modification, some instrument makers have placed bevels in various portions of the instrument, such as where the soundboard of the instrument joins the sides. Under the traditional design, the soundboard and sides of the instrument are at a right angle to each other, defining a common edge. This common edge may have a binding material, such a wood, plastic, or other trim, which protects the soundboard and side from impact damage. Some instrument makers “soften” this common edge by forming a rounded transition section from the side to the soundboard, forming a beveled edge (i.e. a bevel) for a portion of the common edge.\nThe inventor herein has invented a structural feature referred to as the “Ryan bevel” which has since been utilized and/or modified by others practicing the guitar making art. The Ryan bevel generally comprises an approximate 45 degree chamfer on the bass side lower bout. The Ryan bevel feathers in just above the bass-side waist and feathers out just beyond the centerline of the guitar at the tailblock area. The Ryan bevel has some acoustic advantages but it also creates greater comfort for the player since the player's right arm is resting on a chamfered, wide surface rather than against the relatively sharp corner of the guitar where the sides and soundboard of a conventional acoustic guitar meet.\nAnother modification made in recent years to acoustic flat-topped guitars has been the placing of an opening or cutout in the bass side upper bout of the guitar, where this opening is referred to as a “soundport”. The soundport, as with the soundhole beneath the strings in the soundboard, provides an emission means for the sound waves generated by the strings as amplified by the sound chamber. For soundports located in the usual position on the bass side upper bout, sound waves emitted from the soundport are directed toward the face of the guitar player.\nThe known soundports present several problems and disadvantages. First, placing a soundport in the side of the instrument can result in diminished structural strength. For example, the sides of a guitar typically have a wall thickness ranging from 0.060″ to 0.090″, with about 0.075″ being the most common thickness. Given this relatively thin wall, and the continuous load imposed on a guitar body by the string tension—approximately 180 pounds for a steel string guitar—the structural integrity of the load-bearing components of the instrument is always a concern. However, when a side has one or more cutouts for the soundports, the structural integrity may be adversely impacted.\nOne method of resolving this problem is to place a reinforcement member or other structure within the interior of the sound chamber, such as thicker backing on the interior side of the bout in which the soundport will be placed. However, this method has its own disadvantages. Utilizing a reinforcement member in this manner has the disadvantage of adding to the weight of the instrument and/or adversely impacting the acoustic properties of the sound chamber. Adding structural reinforcement on the inside of the instrument utilizes additional materials and is time consuming, thereby resulting in greater manufacturing expense. A reinforcement member can also create additional stresses on the side of the instrument, because the additional piece of wood glued to the side will not always shrink and expand at the same rate as the side material, which can create stresses and possible cracks in the side during relative humidity fluctuations.\nAnother problem presented by the known soundports is that when the soundports are located in the upper bout of the instrument, they are located in a portion of the instrument which is not as acoustically active as other portions of the sound chamber, such as in the lower bout behind and/or adjacent to the bridge.\nThe orientation of the known soundports is typically directed upward into the face of the player, with the axis of the soundport generally perpendicular to the front or soundboard of the instrument. It is to be appreciated that because an audience is typically located at the front of the instrument, the axis of a soundport located in the upper bout of the instrument will be perpendicular to the general direction of the audience. Moreover, depending upon the position in which the player holds the instrument, this location is not optimal for directing sound waves to the player's ear."}
{"text": "The present invention relates to a valve mechanism, particularly to a valve mechanism which can be used for a tube-type fluid container.\nAs this type of valve mechanism, for example, as described in Japanese Patent Laid-open No. 2001-179139, a valve mechanism having a spherical valve body and a spring for giving momentum to the valve body toward a valve seat has been used. Manufacturing costs of the valve mechanism using the spherical valve body and the spring, however, tend to be high.\nConsequently, a valve mechanism having a resinous valve seat, and a resinous valve body which moves between a closed position in which the valve body contacts the valve seat and an open position in which the valve body separates from the valve seat is commonly used.\nIn the resinous valve mechanism, it is preferred that the valve mechanism has a simple configuration which can close a fluid flow reliably. Additionally, it is preferred that the configuration can alter a flow rate of the fluid passing through the valve mechanism discretionally according to a pressure applied to the fluid. As matters stand, however, a valve mechanism satisfying these requirements is not reported."}
{"text": "1. Field of the Invention\nThe present invention is related to an electromagnetic valve apparatus for driving an open/close valve to open or close a high-pressure fuel passage in order to control fuel injection timing and fuel injection quantity used particularly in the fuel injection equipment of an internal combustion engine.\n2. Description of the Related Art\nElectromagnetic valve apparatus has been known for driving by utilizing electromagnetic force, an open/close valve to open or close a high-pressure fuel passage in order to control fuel injection timing and fuel injection quantity. The electromagnetic valve apparatus comprises a core with a coil embedded therein and an armature located adjacent to the core. The armature is attracted toward the core when magnetic flux generated by electric current in the coil passes through the armature, and thus a valve body connected to the armature is driven.\nAs the electromagnetic valve apparatus is composed such that the valve body sits on the valve seat in a state in which the armature is attracted toward the core to close the fuel passage, there must be a gap between the core and the armature in the state in which the armature is attracted toward the core. It is required that this gap is as small as possible for securing strong electromagnetic attraction.\nHowever, heretofore, both the surfaces of the core and the armature facing each other have been shaped as flat surfaces (see FIGS. 1, 3, 4, and 5 of Japanese Laid-Open Patent Application No. 2003-193939, and FIGS. 1, 2, and 3 of Japanese Laid-Open Patent Application No. 2002-98024, for example). The valve apparatus has been composed to provide a sufficient initial gap between the surfaces of the core and the armature facing each other so that the end surface of the armature does not interfere with the surface of the core facing the end surface of the armature even when the armature is installed aslant. This is caused by insufficient accuracy of the component parts, since it is difficult from the standpoint of accuracy of the component parts to compose them so that the gap is even allover the surfaces."}
{"text": "This invention relates generally to thin flexible plastic bags, and more particularly to a plastic bag which can be evacuated after articles are inserted therein.\nHeretofore, the evacuation of entrapped air from plastic bags, after articles have been inserted therein, has been accomplished with the use of complex and expensive machinery. For example, plastic bags have been collapsed about articles by inserting a needle-like tube into the bag and connecting the tube to a vacuum pump for evacuating the bag. This type of machinery is common in the packaging of food products such as cheese, meats and vegetables.\nThe expensive machinery required to evacuate bags heretofore utilized in the packaging of products is not readily available to the individual consumer who, at one time or another, may want to store food products either in a refrigerator or in a freezer by using an evacuated bag."}
{"text": "1. Field of the Invention\nThe present invention relates to a dimming method of an organic EL (electroluminescence) displaying apparatus as one kind of restoring techniques of the organic EL displaying apparatus in which organic EL elements are arranged in a matrix form.\n2. Description of the Related Art\nAn organic EL displaying apparatus is an apparatus constructed in such a manner that a plurality of pixels each having an organic EL element and a driving circuit for driving the organic EL element are arranged on a substrate in a matrix form.\nThere is a case where a defect occurs in the driving circuit provided every pixel due to a defect in a manufacturing step of the organic EL displaying apparatus, so that defective display occurs in the organic EL displaying apparatus. As an example of the defective display, there is such a defect (defective bright-spot) that a bright-spot pixel in which the organic EL element is turned on even in a black displaying mode occurs.\nJapanese Patent Application Laid-Open No. 2004-006339 discloses a method whereby, in an organic EL displaying apparatus, a laser beam is selectively irradiated to a part of a semiconductor layer included in a thin film transistor (hereinbelow, called a TFT) in a driving circuit which a bright-spot pixel has. By irradiating the laser beam to (a part of) the semiconductor layer, a resistance of the laser-irradiated semiconductor layer increases, so that dimming of the bright-spot pixel can be performed (that is, light of the bright-spot pixel can be dimmed).\nFurther, Japanese Patent Application Laid-Open No. 2004-006339 also discloses a method whereby a region to which the laser beam is irradiated is set to a region where a gate electrode layer is arranged upward in the semiconductor layer included in the TFT. According to such a method, the laser beam irradiated from a downward direction is reflected by the gate electrode layer, thereby enabling the dimming to be efficiently performed (that is, enabling the light to be efficiently dimmed).\nHowever, if the method disclosed in Japanese Patent Application Laid-Open No. 2004-006339 is applied to the organic EL displaying apparatus, there are the following problems. That is, when the laser beam is irradiated from the downward direction of the substrate toward the semiconductor layer in order to perform the dimming of the bright-spot pixel, there is a possibility that the laser beam is irradiated to the region where the gate electrode layer provided over the semiconductor layer is arranged. There is such a problem that if the laser beam is irradiated to the gate electrode layer, the semiconductor layer and the gate electrode layer are short-circuited, so that the dimming of the bright-spot pixel is not performed.\nThere is also a case where even if the resistance of the TFT semiconductor layer is increased, the dimming of the bright-spot pixel is not sufficiently performed due to a variation in shape of the TFT by a process variation, a variation in characteristics, or a variation in laser beam which is irradiated. Therefore, it is necessary to increase an energy of the laser beam which is irradiated. In this case, according to the method in the related art, that is, the method whereby the laser beam is irradiated from the downward direction of the substrate toward the region where the gate electrode layer provided over the semiconductor layer is arranged, the foregoing problem of the short-circuit between the semiconductor layer and the gate electrode layer becomes more conspicuous."}
{"text": "1. Field of the Invention\nThe present invention relates to a superconducting X-ray detection apparatus having a highly sensitive superconducting X-ray detector produced using a superconductive material, and to a method for mounting a low-temperature first-stage amplifier for amplifying a feeble signal output from the detector. The present invention also relates to a superconducting X-ray analyzer equipped with the superconducting X-ray detection apparatus.\n2. Background Art\nRecently, superconducting X-ray detectors called TES (Transition Edge Sensor) or STJ (Superconducting Tunnel Junction) have been developed. Currently, study on systemization of the superconducting X-ray detectors is made for meeting various requirements because of their high sensitivity and high energy resolution compared with conventional X-ray detectors. Since each of the superconducting X-ray detectors is a device utilizing superconduction, it needs to be cooled to a low temperature to be driven. Since drive at a lower temperature provides reduction in influence of thermal noise and expectable improvement in performance, the detectors are typically cooled to a very low temperature of about 100 mK.\nTES is a temperature sensor utilizing change of a resistance value occurred when a material in a superconducting state transits into a normal conduction state following temperature rise, and is an X-ray sensor that converts energy of X-ray irradiation into heat and thus performs the detection. When a working point is fixed to a superconducting transition edge by applying appropriate bias current, abrupt temperature rise occurs when X-ray is irradiated on the TES, resulting in change of resistance. Since a current value flowing into the TES changes due to the resistance change, X-ray can be detected by monitoring the current value. Since the change of the current following the X-ray irradiation is feeble, several picoamperes, a SQUID array is typically used for signal amplification.\nAs a main application of the superconducting X-ray detector, an X-ray analyzer and a cosmic-ray detection apparatus are given. When used in the X-ray analyzer, for example, the superconducting X-ray detector is used in combination with a scanning electron microscopy (SEM). The superconducting X-ray detector is inserted into a sample room of the SEM, which is maintained in vacuum environment, at a normal temperature.\nThe detector needs to be placed near a sample to improve X-ray detection efficiency. X-ray generated from the sample by an electron beam irradiation is detected by the superconducting X-ray detector.\nNear the sample, in addition to a sample stage for placing the sample, a lens barrel comprising an electron gun for irradiating an electron beam and an electron lens for focusing the electron beam, and a secondary electron detector for detecting electrons generated from the sample are placed, therefore a space for placing a refrigerator for cooling the superconducting X-ray detector to 100 mK or lower can not be secured. Therefore, the refrigerator is provided separately from the SEM, a cooling head is provided in a thin cylinder called snout constructed in a heat insulation structure, and the superconducting X-ray detector is fixed to an end of the cooling head. The sample is placed near the superconducting X-ray analyzer by inserting the snout from a port of the sample room of the SEM.\nFIG. 5 shows a structure of an end portion of a conventional superconducting X-ray detection apparatus (for example, see non-patent literature 1). The end portion mainly comprises an outer layer of a snout 102 for maintaining vacuum, a radiation-heat shield plate 70 in about two layers for shielding radiation heat, a cooling head 103 for mounting the detector, and a super-insulation 80 applied to inner and outer layers of the radiation-heat shield plate 70. The outer layer of the snout 102 and the radiation-heat shield plate 70 have a certain thickness, and an appropriate interval is necessary, therefore even if the snout has a diameter of about 30 mm, diameter of the cooling head is small, about 10 mm.\nA superconducting X-ray detector 10 such as TES is adhered to a sensor holder 30 produced using a material having high heat conductivity such as oxygen free copper or sapphire together with a connection pad 40, and mounted on an end section portion of the cooling head 30 by a screw clamp.\nA low-temperature first-stage amplifier 20 is adhered to an amplifier holder 21 produced using the same material as the sensor holder 30 together with a connection pad 41, and mounted on the side near the end of the cooling head 30.\nThe superconducting X-ray detector 10 is connected to the connection pad 40 by wire bonding. Among electrodes of the connection pad 40, an output terminal of the superconducting X-ray detector 10 is connected to the connection pad 41 on the amplifier holder 21 using a lead, and as wiring for driving the superconducting X-ray detector 10, connection is made to a connector attached to outer packaging of the refrigerator using a lead.\nOn the other hand, the low-temperature first-stage amplifier 20 is connected to the connection pad 41 by the wire bonding. Among electrodes of the connection pad 41, an electrode for driving the low-temperature first-stage amplifier 20 is connected to a connector attached to the outer packaging of the refrigerator using a lead.\n[Non-Patent Literature 1]\nY. Ishisaki, et al. “Proceeding of the Society of Photo-Optical Instrumentation Engineers (SPIE)”, Hawaii U.S., 2002, p831.\nIn the conventional superconducting X-ray detection apparatus, there has been a problem that since there is no space except for the space for mounting the superconducting X-ray detector, the low-temperature first-stage amplifier that amplifies an output signal from the detector is placed distantly from the detector, and the detector is connected to the amplifier using the lead having high electric conductivity, therefore wiring length is large, the output from the detector that is the feeble signal is influenced by disturbance noise. Furthermore, there has been a problem that handling performance is bad because connection of wiring is complicated.\nThe present invention overcomes the problems in the conventional art by providing a superconducting X-ray detection apparatus having a structure where the influence of the disturbance noise on the output from the detector is minimized by placing the detector and the amplifier in a same place and by providing a mounting method in which the connection of wiring can be easily made."}
{"text": "1. Field of the Invention\nThe present invention relates to a droplet generating method and apparatus therefor, and more particularly to a method and apparatus therefor for emitting a droplet from an orifice of a small liquid chamber containing liquid therein more specifically the present invention provides a method for repeated high-speed generation of the droplets and apparatus for allowing accurate generation of droplets of a uniform diameter.\n2. Description of the Prior Art\nIn the related field there are already known ink jet recording apparatus for example of drop-on-demand type, which are recently attracting particular attention because of their negligibly low noise and absence of an unnecessary ink deposit when recording. Such recording is considered particularly useful in the ability of recording on plain paper without a particular fixing treatment. In the field of such drop-on-demand type ink jet recording there have been proposed various apparatus some of which are already in commercial use, while others are still in the course of development.\nIn summary, the ink jet recording of the drop-on-demand type performs recording by emitting a droplet of recording liquid, called ink, from a small orifice in response to an instruction signal and depositing the droplet onto a recording material. In the known methods, the droplet generation is achieved for example by the use of a piezoelectric element.\nThe present invention relates to a novel drop-on-demand type ink jet recording method which is different from the conventional method utilizing the piezoelectric element and which effects the droplet emission from a small orifice by applying a drive signal to the liquid introduced to a small liquid chamber thereby causing bubble formation in the liquid. Also in a related field U.S. Pat. No. 3,878,519 discloses another apparatus, which while not requiring pressurizing and deflecting means, does provide droplets of insufficient evenness because of a weak droplet forming force and further requires a liquid recovery mechanism for unused droplets, and thus cannot be compactized."}
{"text": "It is a matter of common knowledge to use a magnetic actuator with high force density to operate moving contacts for a purpose of electrical power interruption in the medium-voltage field of technology. Known magnetic actuators have a design with a fixed core in the center of the device, and two moveable plungers, one above and one below the core, that are connected with a plunger rod. Such a device is supposed to generate a high static holding force in the closed position to latch opening and contact springs. The magnitude of this static holding force is the key parameter for the design of the entire circuit breakers and for space and weights reasons it is generally advantageous to generate this force with a small magnetic actuator. In the open position, a lower static holding force is needed to keep the circuit breaker in open position. For bringing the actuator from close to open position feeding the electrical coil of the actuator with electrical energy is needed.\nThe document EP 0 898 780 B1 describes a magnetic actuator with a ferromagnetic armature which is displaceable linearly between two limit positions and which is mechanically connected to a circuit breaker and which in the limit positions is under the influence of magnetically generated forces. The armature and the ferromagnetic shunt body are arranged in succession in a space between first and second abutment. The abutments are pole surfaces of magnetic circuits which include at least one permanent magnet for generating a holding force for the armature. This known device is as well supposed to operate a vacuum circuit breaker. In the closed position, the ferromagnetic shunt body is apart from the armature. The shunt can now be moved towards the armature to initiate the opening operation of the circuit breaker. The known solution is based on a design that does not use the full potential of the static holding force as the effective area between the moveable armature and the fixed yoke is limited to the area that is inside the coil. As a result, the actuator is almost twice as big as needed.\nWO 03/030188 A1 discloses a further magnetic actuator, such as for a vacuum circuit breaker design being large in size. Two electrical coils are needed in order to operate the magnetic actuator or bringing a connected circuit breaker from an open to a closed switching position. A first magnetic flux is generated by the armature and the yoke in such a way that the armature is held in one limit position and the electrical coil generates a second magnetic flux that actuates the armature. The permanent magnet is located between the yoke and a fixed magnetic return element, in such a way that the magnetic flux runs via the magnetic return element. In addition, the armature outside the yoke covers a front face of the yoke, said face running perpendiculary to the direction of displacement of the armature. Since the permanent magnet is provided to hold the magnetic armature in one of the two limit positions, neither mechanically latching nor a constant electrical current supply is specified.\nAlso this known solution uses the armature for generating the static holding force in both limit positions. This solution implies a second magnetic path from the magnets to the armature that is only effective in the open limit position. This second magnetic path increases sizes again and weights of the magnetic actuator. The known solution also specifies a closed room around both the armature. The ferromagnetic shunt body forms the two abutments that need to fulfill magnetic functions. This increases the size and weight of the actuator further. The known solution entails the driving of the ferromagnetic shunt body back to the lower abutment during the opening operation. This driving should include additional energy that is not available for the opening operation, which is the most critical operation of a circuit breaker in case of short circuit switching."}
{"text": "Image display devices such as cathode ray tubes (CRTs), liquid crystal display devices (LCDs), plasma display devices (PDPs), electroluminescence display devices (ELDs), and electronic paper are generally provided with optical laminated bodies for antireflection on the outermost surface thereof. Such optical laminated bodies for antireflection suppress reflection of images and decrease the reflectivity by light diffusion or interference of light.\nAs one of such antireflection optical laminated bodies, an anti-glare film including an anti-glare layer that has a surface roughness and is provided on a surface of a transparent substrate is known. Such a surface roughness of an anti-glare film diffuses ambient light, whereby a decrease in visibility can be prevented.\nConventional anti-glare films are known in which anti-glare layers formed by applying resins containing filler such as silicon dioxide (silica) on surfaces of transparent base films (see, for example, Patent Literatures 1 and 2).\nFor example, in such anti-glare films, a surface roughness is formed by adding cohesive particles, an inorganic filler and/or an organic filler in a resin; laminating a film with unevenness on an anti-glare layer; and phase separation using compatibility among compounds such as two or more different polymers, which constitute a binder.\nIn all such conventional anti-glare films, a light diffusion function and an anti-glare function are obtained by the effect of a surface shape of the anti-glare layer. In order to increase such an anti-glare function, the size of the surface roughness needs to be increased. However, the large-sized projection and depression increases a haze of a coating, whereby white muddiness is generated.\nFurther, the twinkling brightness which is called scintillation is generated on the surfaces of conventional anti-glare films, which results in a decrease in visibility of display screens.\nIn recent years, high definition liquid crystal displays have been developed. However, scintillation is generated in such high definition liquid crystal displays when conventional anti-glare films are used therein. Such scintillation is a serious problem.\nFurthermore, in the conventional anti-glare films, cracks are generated during a process for adhering the films to polarizing elements or a fabrication process of liquid crystal cells, at, for example, an interface between light-transmitting fine particles and a light-transmitting resin that are included in an anti-glare layer. Furthermore, in the conventional anti-glare films, anti-glare performance and a scintillation state are changed over time due to the haze variation in response to temperature and moisture changes. Therefore, the conventional anti-glare films are less resistant to moisture and heat.\nPatent Literature 3 discloses an anti-glare material obtainable by mixing a binder resin with resin beads in which 70% or more thereof is swollen with a solvent.\nIn an anti-glare film provided with an anti-glare layer including such resin beads that are previously swollen with a solvent, the adhesion at the interface between the resin beads and the binder resin is expected to be improved. Such an anti-glare film is expected to be used in high definition displays.\nIn such an anti-glare film provided with an anti-glare layer including resin beads that are previously swollen with a solvent, the adhesion at the interface between the swollen resin beads and the binder resin in the anti-glare layer is improved only by an anchor effect created at the interface. Therefore, the adhesion and the like can be further improved.\nPatent Literature 4 discloses a method for producing a light diffusion film in which the density of light-transmitting fine particles, a coating composition, and the average particle size of light-transmitting fine particles are in a specific relation to one another. Patent Literature 5 discloses an optical film including light-transmitting particles having a certain average particle size that depends on the thickness of a light diffusion layer. Patent Literature 6 discloses an optical film including an anti-glare layer formed from a coating composition in which a binder and two types of resin fine particles different in particle sizes and swelling rates that satisfy a certain relation with each other are dispersed in a dispersion. Patent Literature 7 discloses a method for producing a light diffusion film, including a coating composition for a light diffusion layer. The coating composition includes light-transmitting fine particles, a light-transmitting resin containing a certain amount of a light-transmitting polymer with a molecular weight of 1000 or more, and a solvent.\nHowever, Patent Literatures 4 to 7 do not disclose at all that light-transmitting fine particles are impregnated with a binder, and do not disclose at all, either the effect exerted by the impregnation of light-transmitting fine particles with a binder. Therefore, the anti-glare films disclosed in the patent literatures do not sufficiently achieve an anti-glare property, prevention of white muddiness, and prevention of scintillation.\nTherefore, in conventional anti-glare films, generation of cracks at the interface between resin beads and a binder resin is required to be prevented at a high level in a process for attaching the films to polarizing elements and a fabrication process of liquid crystal cells. Furthermore, white muddiness is required to be reduced by reducing reflection at the interface between a binder and resin beads and scintillation is also required to be prevented by appropriately dispersing resin beads."}
{"text": "Computing devices can be used to request content from other computing devices over a communication network. In a common application, a client computing device can request a web page from a server computing device via the internet. Browser application software on the client computing device typically retrieves a requested web page, processes resource identifiers embedded in the web page to generate requests for additional resources (e.g., images, script files, etc.), and renders the content for presentation. From the perspective of a user of a client computing device, a user experience can be defined in terms of the performance and latencies associated with obtaining and rendering the requested network content on the client computing device. Latencies and performance limitations of any of the above processes may diminish the user experience.\nOptimizations and other improvements may be implemented to reduce latency and otherwise improve the user experience. For example, some content items may be cached at a client device, and future requests for the content items can be fulfilled from the local cache. As another example, proxy servers may be used to cache content for multiple clients and then pass the content to the clients that originally requested the content. Subsequent requests to the proxy server can be served from the proxy cache if the content is present, rather than requiring retrieval of the content from an origin server. As a further example, content delivery networks (“CDNs”) can locate points-of-presence closer (in either a geographic or networking context) to certain client devices. Content requests may be made to the CDN rather than the origin server, and the CDN can fulfill the request more quickly due to being closer to the requesting device. Some proxy servers and other intermediary systems can pre-render or otherwise pre-process content prior to responding to a client request. Such pre-processing enables some of the processing burden to be offloaded from the client device to the intermediary system."}
{"text": "This specification relates to communications between virtual machines.\nCloud computing is network-based computing in which collections of servers housed in data centers or “server farms” provide computational resources and data storage as needed to remote end users. Some cloud computing services provide access to software applications such as word processors and other commonly used applications to end users who interface with the applications through web browsers or other client-side software. Users' electronic data files are usually stored in the server farm rather than on the users' computing devices. Maintaining software applications and user data on a server farm simplifies management of end user computing devices. Some cloud computing services allow end users to execute software applications in virtual machines. In a public cloud computing environment, multiple users are able to launch virtual machines (VMs), and each VM launched by a user is included in a cluster of other VMs launched by the user.\nLike reference numbers and designations in the various drawings indicate like elements."}
{"text": "1. Field of the Invention\nThe invention relates to a method for determining a weight per unit area and/or a chemical composition of a conveyed material sample, and to a device for carrying out the method.\n2. Discussion of the Background\nFor determining a weight per unit area (e.g., grams per square meter; g/m2) and/or a chemical composition of material samples or material webs in conveyor and production plants partially operated at high speeds, absorption measurements are carried out according to the state of the art, using ionized radiation, in particular x-ray, gamma or particulate radiation. In this case, the webs or samples to be tested are passed through between a source of radiation and a detector arrangement placed opposite the source of radiation, and, in a transmission measurement, the intensity of the radiation absorption is determined. This radiation absorption depends on the weight per unit area of the webs and samples so that the mentioned characteristics of the webs and samples may be measured in a contact-free manner by means of this known method and corresponding devices.\nA further possibility of measuring the weight per unit area of a material web or a sheet material consists in detecting scattered radiation. NDC Infrared Engineering and Adaptive Technologies, for example, are companies who use this method in their products by detecting back-scattered photons originating from a gamma radiation source, e.g. photons of the 59 keV line of the Am 241 isotope.\nA systematic measurement error occurs in this method when the distance between the place where the particles of the incident ray are scattered and the location of the detector varies. Such a distance variation may occur due to a random fluttering movement of the material web, and also always occurs when the thickness of the measured product changes. In this case, purely geometric reasons cause the measurement signal to be influenced, and these reasons, if not corrected by additional independent measurements, will lead to an erroneous reproduction of the value of the weight per unit area. This systematic effect, however, may be limited when a large distance between the measured product and the detector is selected. But the distance may not be so large, for the purpose of maximizing the detection efficiency of the scattered radiation and lowering the associated statistical error of the detector with using, at the same time, the least possible intensity of the source of the incident ray. It is therefore important to correct the inaccuracies caused by distance variations in order to attain minimal measurement error at an economic cost.\nIn all known methods for measuring weights per unit area in which ionized radiation is used, the radiation originates from the decomposition of radioactive isotopes. In this case, the operating company is subjected to the radiation-protection regulation requiring the operating company to observe rest times even during a standstill of the plant, as well as being subject to elaborate safety measures.\nThe economic availability and the natural properties of radioactive nuclides and the ionized radiation thereof, e.g. half-life, type and energy distribution, strongly restrict the number of industrially usable preparations. In addition, using a single radionuclide may allow only a restricted functional range to be covered. Thus, the beta radiation of the frequently used nuclide Kr 85, for example, having an end-point energy of 687.4 keV, due to its absorption behaviour is well suited for determining weights per unit area of up to about 1000 g/m2. Higher weights per unit area may no longer be detected with it in a satisfactory manner and are preferably measured with the nuclide Sr 90, in which in the consecutive decomposition of Y 90, a beta radiation of a higher energy and hence of a higher penetration occurs having an end-point energy of 2280.1 keV."}
{"text": "The present invention relates to an exercise assembly, and more particularly, to a modular exercise assembly intended for use in water.\nThe past few years have seen an increase in the awareness of individuals of their physical well-being and stamina. The benefits of exercise have been highly publicized and the exercise market has increased tremendously as a result of this publicity and awareness. Exercise may take a variety of forms, and includes cardiovascular exercise such as jogging, running or cycling, or it may include traditional weight or resistance training such as weightlifting. Whereas cardiovascular exercise typically improves endurance, weight or resistance training builds up or strengthens particular aspects of an individual's body. The present invention is primarily directed to this latter category of exercise.\nWeight training usually requires a barbell shaft and a plurality of different weighted elements which are held on the ends of the barbell shaft by retaining collars or the like. The user lifts, or otherwise moves, the barbell between at least two preselected positions for a predetermined number of repetitions. A user may perform a large number of repetitions at a moderate weight to tone specific body areas, or the user may perform a number of repetitions with maximum weight until failure in order to build up strength in his particular muscles.\nResistance training is similar to weight training in that it requires a user to perform repetitions. However, instead of using static weights, the user typically pushes or pulls against a resistance force of an exercise machine which may be supplied by rubber or elastic bands. Although weight and resistance training are generally effective methods of physical conditioning, they are not without certain disadvantages. For example, weight training necessitates that a user keep a large number of weights on hand. The weights are bulky and heavy in nature and are not easily stored or transported. Additionally, most weight training exercises require the presence of a spotter, who is another individual who is available to assist the user in case of trouble. Also, both resistance and weight training tend to impart, dependent on the particular exercise, undue stress on the user's joints. Resistance training not only requires a machine, but also necessitates replacement of the resistance bands as they wear down over time.\nHydrodynamic exercises, where the user is partially immersed in water and uses the water as a weight or resistance medium, has been recognized as providing comparable benefits to weight or resistance training while eliminating the detrimental stress on the joints of the user. Hydrodynamic exercise includes either the user lifting or pushing one of his limbs, such as an arm or leg against the water in a repetitive movement or the user manipulating an exercise device in the water in a repetitive manner. These hydrodynamic exercise devices are well known in the art and include, for example, barbells having one or more foam rubber elements which the user pushes or pulls against the resistance of the water as disclosed in U.S. Pat. No. 5,188,993, issued Feb. 9, 1993. Hydrodynamic exercise devices may also include a barbell or dumbbell having plurality of planar fin-or blade-like members which the user pushes and/or pulls against the water, such as those disclosed in either U.S. Pat. No. 4,623,142, issued Nov. 18, 1986 or U.S. Pat. No. 4,819,915, issued Apr. 11, 1989.\nThese types of hydrodynamic exercise equipment are easy to use and relatively simple in design and construction. However, the blades or fins which provide the water resistance surface(s) for the exercise device typically extend radially outwardly from a shaft at different angles. These extensions render the exercise equipment bulky and cumbersome for transport and storage. Heretofore, no one has provided aquatic exercise equipment which is modular in nature which may be easily disassembled for component storage and transport, and which may be assembled in a minimum amount of time.\nIt is therefore an object of the present invention to provide a modular aquatic exercise equipment assembly in which the exercise equipment assembly may be easily disassembled and which requires a minimum of space for storage and transport.\nAnother object of the present invention is to provide an improved hydrodynamic exercise equipment component having a user-engageable shaft and two water-resistance members removeably attached thereto at opposite ends of the shaft, each of the water-resistance members including planar fins which radially extend outwardly therefrom in different radial directions, and each of the water-resistance members including a hand grip portion interposed between the fins.\nYet another object of the present invention is to provide a set of components for a modular assembly of aquatic exercise equipment which may be quickly assembled together to form an aquatic exercise device and which may be quickly disassembled into its constituent components for storage or transport space, the disassembled components occupying a minimum of space, the set of components including two aquatic dumbbell members threadedly engaging a barbell shaft, each of the aquatic dumbbell members having a hand grip portion thereon provided near the intersection of two engageable half portions, each half portion having a planar surface constituting a fin extending radially outwardly therefrom, the two half portions being mutually engageable with each other such that, when engaged together, the half portions define a dumbbell having a plurality of fins extending radially outwardly from the hand grip portion, the fins being displayed at different angles to each other.\nStill yet another object of the present invention is to provide hydrodynamic exercise equipment in which two water engagement members are removeably attached to opposite ends of the assembly and in which the water engagement members include two interlocking members which engage each other along an axis of the assembly, each of the interlocking members including a planar fin portion and an elongated shaft portion disposed centrally therein, the shaft engagement portions engaging each other such that the planar fin portions extend radially outwardly therefrom along different radial directions to define a plurality of water-resistance surfaces of the assembly."}
{"text": "There has already been proposed a device of the type generally described above, featuring a cylindrical tube inclined from the vertical below a sliver guide, the bottom end of the tube being mounted on a supporting element, and the top end or upper part being accommodated in a ball bearing, the outer race of which is located in a horizontal plate describing a gyratory movement through the action of two crankshafts coupled by drive elements, such that the tube will undergo a generally nutational movement in operation, in the sense that the tube will move in a generally conical path about its lower end. An exemplary such device is disclosed in a commonly assigned U.S. application Ser. No. 88,814 filed Oct. 29, 1979, (Schopper), entitled \"Apparatus for loading a rope-like strand of fibers into spinning cans.\" Attention is also invited to commonly assigned earlier U.S. application Ser. No. 920,822, filed June 30, 1978, now U.S. Pat. No. 4,173,057, corresponding to Swiss Pat. No. 611,239. With such a device in operation, the sliver is piled into a can in cycloidal loops, but there is a tendency for the sliver to accumulate or concentrate in particular areas of the can in an undesired manner, the sliver typically accumulating in the center of the can. This applies in particular to the situation in which the tilted tube and the cans at the delivery end of the device are of relatively small diameters, as in feeding open-end frames.\nAn object of the instant invention is to remove this shortcoming in the piling operation, thereby bringing about improved distribution of the sliver in the piling operation, principally by shifting the axis of the cycloidal loops. In general, this object is met by locating the crankshafts on a second horizontal plate which also describes a gyratory movement through a second pair of coupled crankshafts, the first pair of crankshafts being driven more slowly than the second pair. Preferably the first pair of crankshafts is rotated between ten degrees and fifty degrees for every complete rotation of the second pair of crankshafts.\nOther and further objects, advantages and features of the instant invention will be apparent to those skilled in the art from the ensuing description of a preferred embodiment, taken in conjunction with the appended drawing."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates to a method of utilizing a mobile device to detect an emergency, and more particularly, to a method of detecting an emergency and determining whether a user is involved in the emergency according to an acceleration status of a mobile device carried by the user.\n2. Description of the Prior Art\nAs the average wealth of a nation increases, people start to place more importance on their quality of life and the security of their family. Security monitoring devices which integrate various fields such as burglar alarms, entrance management and remote monitoring can be utilized for personal safety and protection of property. As the related technology matures, various digitized, networking or integrated security products and applications have become commercially available. Most modern security monitoring products and applications are realized by monitors or security systems. Coverage of such monitors and security systems, however, may not be satisfactory. If an accident or emergency occurs in a place not covered by the monitors or security systems, the circumstances will not be known, and follow-up assistance or appropriate handling cannot be performed. Thus, how to immediately detect and determine whether a user is involved in an emergency has become a main focus of the industry. Such functionality can enable assistance or processing to be performed immediately after an emergency, in order to achieve effective monitoring for personal safety."}
{"text": "As a radiation detector for use in a CT sensor, etc., there is a detector of an arrangement wherein a scintillator is disposed on a light-incident surface of a semiconductor photodetecting element, which is a semiconductor element. With such a radiation detector, when an X-ray, γ-ray, charged particle beam or other radiation to be detected is made incident on the scintillator, scintillation light is generated inside the scintillator by the radiation. The semiconductor photodetecting element then detects the scintillation light that is made incident from the scintillator and output an electrical signal that is in accordance with the intensity of the radiation.\nAlso, in order to carry out signal processing of the electrical signal output from the semiconductor photodetecting element, a signal processing element is provided. As an arrangement for electrically connecting the semiconductor element with the signal processing element and transmitting the electrical signal, an arrangement is used wherein a semiconductor device is arranged in one united body by connecting the semiconductor element to a wiring substrate that is provided with a conduction path and the signal processing element is connected to the wiring substrate of this semiconductor device. Such a semiconductor device using a wiring substrate is used in various applications besides radiation detectors (see for example, Japanese Patent No. 2555720 and Japanese Patent Application Laid-Open No. H3-203341)."}
{"text": "The present invention relates to semiconductor structures, and more particularly to field effect transistor structures having an asymmetric SiGe channel, and methods of manufacturing the same.\nA SiGe channel, i.e., a channel composed of a silicon-germanium alloy, provides a smaller band gap than a channel composed of silicon in a metal-oxide-semiconductor field effect transistor (MOSFET). The SiGe channel can be advantageously employed, for example, in combination with a high-k gate dielectric and a metal gate electrode thereupon to provide an optimized band gap in field effect transistors. However, the reduction in the band gap induces gate-induced drain leakage (GIDL) current in field effect transistors employing a SiGe channel. Such increase in the GIDL current has been observed, for example, in Dongyun Kim et al., “Band to Band Tunneling Study in High Mobility Materials: III-V, Si, Ge and strained SiGe,” Device Research Conference, pp. 57-58 (2007), Krishna C. Saraswat et al., “High Mobility Materials and Novel Device Structures for High Performance Nanoscale MOSFETs,” IDEM (2007), and X. Chen et al., “A cost Effective 32 nm High-k/Metal Gate CMOS Technology for Low Power Applications with Single-Metal/Gate-First Processes, IEEE VLSI pp. 88-89.\nSuch increase in the GIDL current due to a SiGe channel has a negative impact on the performance of a field effect transistor by increasing off-current of the field effect transistor, thereby rendering the field effect transistor unsuitable for low power applications. In order to fully utilize the advantage of a SiGe channel in providing a well controlled threshold voltage, the GIDL current of the transistor needs to be controlled to a minimal level."}
{"text": "As a method of surface treatment of aluminum and/or aluminum alloy, there is known chromate treatment, and zirconium phosphate treatment (cf. Japanese Kokai Publication Hei-05-302042, for instance). However, since chromate treatment uses harmful chromium metal, it causes a problem of liquid waste treatment. In zirconium phosphate treatment, there is a problem of being low in corrosion resistance.\nFurther, though surface treating agents which use a fluorine and zirconium-containing compound are used, there may be cases where the corrosion resistance is inadequate in some uses even though using such surface treating agents. Therefore, a rust prevention coating agent is desired, which can impart higher corrosion resistance to aluminum and/or aluminum alloy. As such a rust prevention coating agent, there has been disclosed a surface treating agent, which contains a zirconium compound and a compound containing nitrogen atoms having a lone pair (cf. Japanese Kokai Publication 2000-204485, for instance). But, adequate corrosion resistance has not been attained even though using such a surface treating agent. And, coats obtained by using such a surface treating agent were inadequate for being adopted as a rust-proofing because they have low corrosion resistance during not yet being coated and insufficient adhesion to a coating film."}
{"text": "Typical water meter assemblies comprise a series of plumbing elements positioned within a meter box. The plumbing includes a water meter coupled to an angle meter stop (sometimes also referred to as an “angle meter valve”). The angle meter stop includes a valve and is typically positioned adjacent to the water meter toward one end of the meter box. The angle meter stop forms a 90° downward bend, with a horizontal portion coupled to the water meter and vertical portion coupled to a downwardly extending vertical conduit that connects with a water main or water source. The vertical portion of the angle meter stop is typically attached to the downwardly extending vertical conduit using a “grip-tite” friction fit.\nPeriodically, maintenance needs to be performed on the plumbing elements of the water meter assembly. For example, the water meter may need to be replaced every 15 years. To replace the water meter, the valve of the angle meter stop is first closed and the old water meter disconnected from the angle meter stop. A new water meter is then installed. With the valve closed, the water pressure in the vertical portion of the valve can exert a strong upward force on the angle meter stop. Without the water meter connection holding the angle meter stop in place, the only force holding the angle meter stop to the vertical conduit is the “grip-tite” friction fit. If the friction fit is not properly formed, or the friction fit fails for any reason, the water pressure can cause the angle meter stop to violently blow off from the vertical conduit and injure a person performing maintenance."}
{"text": "The present invention relates to an electric discharge machining power source of an electric discharge machine which performs machining by using electric discharge of a capacitor.\nA power source circuit of an electric discharge machine which performs machining by using electric discharge of a capacitor conventionally uses a circuit as shown in FIG. 5. Referring to FIG. 5, reference symbol E denotes a power source; R1, a current-limiting resistor; C2, a charge/discharge capacitor; P, an electrode; W, a work; T2, a transistor serving as a switching element; and G2, its base. With this circuit, a pulse is input to the base G2 of the transistor T2 to turn it on and thus to charge the capacitor C2. The charge voltage of the capacitor C2 is applied between the electrode P and the work W. The discharge current from the capacitor C2 flows as an electric discharge between the electrode P and the work W, thereby performing electric discharge machining. In order to quickly charge the capacitor C2, the resistance of the resistor R1 must be set small. However, when the resistance of the resistor R1 is set small, the resistor R1 generates heat, resulting in a large energy loss.\nIn view of this situation, an electric discharge machining power source as shown in FIG. 6 is developed (Japanese Patent Application No. 59-35984). Referring to FIG. 6, reference symbol E denotes a power source; T3 and T4, charge and discharge transistors serving as switching elements; G3 and G4, bases thereof; C3, a charge/discharge capacitor; D3, a diode; P, an electrode; W, a work; and L3, an inductance.\nThe operation of this electric discharge machining power source circuit will be described. First, when a pulse is applied to the base G3 of the charge transistor T3 to turn it on, a current flows from the power source E to the capacitor C3, inductance L3, transistor T3, and power source E, and the capacitor C3 starts to be charged. At this time, the current linearly increases because of the impedance of the stray inductance L3. Then, the charge voltage of the capacitor C3 also gradually increases. When application of the pulse to the base G3 of the charge transistor T3 is stopped to turn it off, the current as the energy accumulated in the stray inductance L3 flows through the diode D3, thus constituting a so-called flywheel circuit to further charge the capacitor C3. In this case, the charge current and voltage to the capacitor C3 are determined by the value of the stray inductance L3 and the width of the pulse applied to the base G3 for turning on the transistor T3, as is apparent from the above description. Therefore, the charge voltage can be adjusted by adjusting the pulse width and the stray inductance L3.\nIn this manner, the capacitor C3 is charged, the pulse is applied to the base G4 of the transistor T4, and thus the transistor T4 is turned on. As a result, the charge voltage in the capacitor C3 is applied in the gap between the work W and the electrode P, the capacitor C3 starts discharging, and the discharge current flows. In this manner, since this electric discharge machining power source uses no power source limiting resistor, no energy loss is caused and good power source efficiency is ensured. However, since the discharge switching element T4 is used, the overall system becomes expensive. Particularly, in order to flow a large discharge current, a plurality of transistors T4 (only one is shown) serving as the switching elements must be connected in parallel with each other, rendering the electric discharge machining power source expensive."}
{"text": "This invention relates to systems for inspection of substrates, especially semiconductor wafers and reticles. More specifically, the invention relates to a novel system which utilizes a sensor array in conjunction with two threshold arrangement, preferably integrated on a CMOS camera.\nSeveral systems are known in the art for the inspection of wafers and reticles. Three examples of such systems that are currently available on the market are depicted in FIGS. 1-3. In the system exemplified in FIG. 1, the wafer 100 is illuminated with a light beam emanating from a light source 110 and reaching the wafer at a 90xc2x0 angle (generally referred to as normal illumination). Preferably, light source 110 provides coherent light, i.e., source 110 may be a laser source. The light beam is scanned over the wafer by a scanner 120, typically an acousto-optical deflector (AOD) or a rotating mirror, in the direction marked by the double-headed arrow. The wafer 100 is moved in the perpendicular direction by moving the stage upon which the wafer rests. Thus, a two dimensional area of the wafer can be scanned by the light beam.\nSince the wafer has basically a mirror-like top surface, the light beam is specularly reflected back per Snell\"\"s law at 180xc2x0. This specularly reflected light is collected by a light sensor 140 and its signal is used to obtain a xe2x80x9cbright fieldxe2x80x9d image, i.e., an image created from specularly reflected light. However, whenever the light beam hits an irregularity on the wafer, such as a particle, the light scatters in various directions. Some of the diffracted/scattered light is then collected by the light sensors 130, and their signal is used to obtain a xe2x80x9cdark fieldxe2x80x9d image, i.e., an image created from diffracted/scattered light. Thus, irregularities appear in the dark field image as stars in a dark sky.\nIn the system exemplified in FIG. 2, the wafer 200 is illuminated by a light beam emanating from light source 210, but reaching the wafer at a shallow angle, generally referred to as grazing illumination. The light beam is scanned over the wafer by a scanner 220, typically an acousto-optical scanner or a rotating mirror, in the direction marked by the double-headed arrow. The wafer 200 is moved in the perpendicular direction by moving the stage upon which the wafer rests. Thus, a two dimensional area of the wafer can be scanned by the light beam.\nSince the light reaches the wafer at a grazing angle xcex8, its specular reflection is at a corresponding angle, xcex8xe2x80x2, according to Snell\"\"s law. This light is collected by sensor 240, and its signal is used to create the bright field image. Any scattered light is collected by sensors 230, the signal of which is used to create dark field images.\nIt should be appreciated that in the above exemplified systems, with respect to each sensor the image data is acquired serially. That is, each two dimensional image, whether bright or dark field, is constructed by acquiring signals of pixel after pixel, per the scanned light beam. This is time consuming serial operation, which directly affects the throughput of such systems. Moreover, the scan speed of such systems is limited by the scanner\"\"s speed (i.e., the band-width for an acousto-optic scanner) and by the electronics that support the detectors, e.g., the PMT (Photo-Multiplier Tube). Thus, a need exists to develop a system that does not utilize a scanned light beam.\nAnother difficulty with systems which use coherent light is diffraction caused by features arranged in a repeated order, thereby effectively forming a grating. Specifically, in semiconductor devices many features are constructed in a repeated order fashion. When these features are illuminated by a coherent light beam, they diffract the light in much the same manner as a diffraction grating would diffract the light. However, such constructive diffraction can be mistaken by the system for a defect. One way to overcome such a problem is to use a spatial filter in the Fourier plane, as exemplified by filters 235 in FIG. 2. This problem and proposed solutions are disclosed in, for example, U.S. Pat. No. 4,898,471, 4,806,774, and 5,276,498, which are incorporated herein by reference.\nThe system depicted in FIG. 3 performs a bright field inspection only, but does not use a scanner to scan the light beam. Instead, light source 310 provides a relatively broader light beam which illuminates the wafer 300 with a relatively large spot 315. A TDI sensor is used to image an elongated strip 325 of the illuminated spot. The length of the strip corresponds to the width of the TDI sensor. For example, if the TDI sensor comprises a 2048xc3x972048 pixels, then the scanned strip is of size 2048xc3x971. As the wafer is moved by the stage, strips are imaged and collected to form a bright field two-dimensional image of the inspected area.\nLooking forward, as design rules shrink, the importance of detecting increasingly small irregularities becomes paramount. With design rules such as 0.18 and 0.15 xcexcm, very small irregularities, such as particles of sub-micron size, can be killer defects and cause the device to malfunction. However, in order to detect such small irregularities, one needs to use a very small wavelength light source, such as ultra violate (UV) or deep ultra violate (DUV) light source. This presents at least two crucial problems: first, optical elements operating in the DUV regime are expensive and, second, small short wave implies small spot size of the light beam; thus, the scanning speed and collection data rate need to be increased.\nThe small size of killer defects also present a formidable challenge for bright field system which do not use scanning, such as the TDI system depicted in FIG. 3. Specifically, since bright field system construct an actual image of the inspected area, the image includes multitude of structural elements built upon the wafer. Thus, the resulting image looks much like a maze, and it is increasingly difficult to detect a small irregularity in the maze-like image. Thus, the system requires a complicated image processing algorithm to recognize the defect, thereby increasing the processing power and time required and increasing the cost of purchasing and operating the system. It is not clear at this time whether even the most sophisticated algorithm may be unable to detect such small irregularities.\nAccording to the present invention, a two dimensional sensor array is used to collect light returned from the inspected substrate. The signal generated by each individual sensor is passed through a threshold. Those signals which are below the threshold are amplified and are summed up. The summed signal is then passed through a second threshold. Summed signals which pass the second threshold are flagged as indicating suspect locations on the substrate.\nThe preferred implementation of the present invention is an integrated CMOS camera. Specifically, both thresholds, the amplifiers and the adder are integrated into the CMOS camera. This implementation is advantageous in that it enables much faster throughput of the system."}
{"text": "The present invention relates generally to an optical device, and more particularly, to an optical device utilizing a superlattice of one or more indirect bandgap materials.\nThe strengths of optical transitions between conduction band states and valence band states in semiconductors are greatly influenced by whether or not these transitions require phonon participation. Direct transitions (i.e. transitions between states with the same wavevector) can occur without phonon participation, and can have relatively large transition strengths. However, indirect transitions (i.e. transitions between states with unequal wave vectors) require phonon emission or phonon absorption in addition to the photon emmision or absorption process, and have relatively weak transition strengths.\nThe efficiency of optical devices can be significantly improved by the use of active region materials with direct transitions at the photon energies of interest. In addition, certain optical devices such as laser diodes essentially require the use of materials with direct band-edge transitions; i.e., electron-hole recombination events which occur at energies close to the band gap.\nMany materials having direct transitions are used in optical devices. For instance, bulk GaAs and InAs are examples of conventional direct transition materials. However, some applications require the use of materials having indirect bandgaps. For example, all of the known technologically relevant compound semiconductors with appropriate bandgaps for green LEDs and green laser diodes have indirect band gaps. The emission efficiencies of sub-bandgap transitions can be enhanced somewhat in certain materials by the incorporation of selected impurities. However, the efficiencies obtained in this manner do not approach those obtained from direct gap materials."}
{"text": "In Australian Patent Number 769719 there is described and claimed a multifocal contact lens wherein the contact lens is made of flexible material and the contact lens is arranged to translate on an eye.\nThe entire disclosure of Australian Patent Number 769719 is incorporated herein by reference.\nThe present invention provides a means of manufacturing a lens of the type described in Australian Patent Number 769719."}
{"text": "Various embodiments of the present invention relate generally to a rifle or firearm and in particular, a gas regulator for a firearm. Still other embodiments relate generally to a rifle or firearm and in particular, a locking sight for a firearm.\nFolding or collapsible sights are desirable for weapons in that the same allows the sight to transition from a first folded or stowed position to a second extended or deployed position. In one application, the first folded or stowed position allows the sight to be configured or located in a position that does not interfere with an optical sighting system or accessory that is mounted to a rail of the weapon. However, folding or collapsible sights may be inadvertently moved from the deployed position. Accordingly, it is desirable to provide a folding or collapsible sight with a locking mechanism that prevents inadvertent movement of the sight\nA flash suppressor and/or a sound suppressor used with a firearm or rifle will increase the rate of fire or cycling of the firearm. This increase in cycling may also increase the wear rate of components of the firearm. In addition, variations in the types of ammunition used with rifle will also vary the cycling rate of the rifle and a rifle configured for use with one type of ammunition may cycle too fast when fired with another type of ammunition thus causing the aforementioned excessive wear to the components of the rifle. Non-limiting examples of different types of ammunition are 5.56 mm NATO ammunition and 223 Remington ammunition each of which may provide different cycling rates when fired from a rifle having the same gas port configuration in a gas impingement system. Still further there are reliability issues associated with higher cycling rates when ammunition is changed for example, round feeding may not occur or may occur improperly.\nAccordingly, it is desirable to provide an apparatus and method for regulating the gas pressure in a rifle and thus allowing the same to be configured for use with suppressors and various types of ammunition."}
{"text": "The present invention relates in general to compressed natural gas (\"CNG\") storage and dispensing systems, and more particularly concerns novel apparatus for dispensing CNG at substantially constant and easily measured mass flow rates.\nIn light of expanding concerns about emissions and foreign energy dependence, commercial and passenger vehicles are increasingly being designed to operate on CNG. There is a growing need, therefore, for CNG fueling stations able to resupply these vehicles. In order for CNG-fueled vehicles to have commercial appeal, not only must these stations be conveniently located, as gasoline filling stations are today, but in addition they must be able to refuel a constant stream of vehicles, each in a relatively brief period of time. Further, they must be reasonably inexpensive to install and operate. Given these considerations and the high pressures at which CNG is typically stored within a vehicle, these CNG fueling stations therefore must be specially designed.\nA paramount concern of the CNG station operator is the station's \"storage effectiveness,\" defined as the fraction of stored CNG that can be transferred, at a particular pressure, to CNG vehicle tanks. Storage effectiveness, together with compressor flow rate characteristics, determine the number of vehicles that a station can fill in a given time period. Today, storage effectiveness can be improved by increasing the CNG storage pressure, the number of \"cascade\" levels, or both.\nIn a typical cascade CNG station, empty vehicle tanks are filled first from one of a series of three conventional CNG storage tanks. If the pressures in the vehicle and first storage tanks equalize at a pressure below the maximum desired vehicle tank pressure, a sequential valve then connects the vehicle tank to a second storage tank, which contains CNG at higher pressure. If necessary, this process then repeats using a third tank. A dome valve ensures that the vehicle tank pressure does not exceed the maximum desired pressure. During filling, a priority valve determines, based on the pressures in each, which storage tank should be refilled first by the compressor. Throughout the filling operation, a mass flow sensor monitors the total amount of CNG transferred to the vehicle.\nAlthough they greatly increase both the cost and the complexity of the station, these various valves and the gas flow rate sensor are essential to the operation of the cascade system. Even with this expense, however, the storage effectiveness still remains well below unity. Storage effectiveness can be incrementally improved in this type of system only by increasing CNG storage pressure. However, not only does this require a higher pressure compressor, but any energy used to increase storage pressure above the maximum vehicle storage pressure is necessarily lost during the filling operation. This increases the total operating cost of filling each vehicle.\nVehicle owners and station operators would also prefer to reduce the temperature rise in the vehicle tank during the filling operation. The lower the final temperature of the tank, the greater the mass of CNG that can be stored there. When filled by a cascade system, however, vehicle tanks experience considerable \"compression heating.\""}
{"text": "Such bowling game systems are configured with: an approach where players carry out a bowling action; a lane extending from the approach and on which the ball bowled by the player rolls; a plurality of pins placed on an end of the lane opposite a side of the approach; and trough-like gutters provided paralleling both sides of the lane to receive the ball slipping off the lane. Players enjoy the game by competing with each other for a higher score that depends on the number of pins they knock down (see Patent document 1).\nPatent document 1: Japanese Unexamined Patent Application Publication No. 2005-334224"}
{"text": "The embodiments described herein relate generally to sampling devices, and more specifically, to a sampling device including a mechanical force feedback mechanism.\nScreening systems for threat compounds such as explosives as well as chemical and biological weapons are commonly used. These screening systems must be able to collect, concentrate, and analyze trace samples quickly and accurately. Many detection technologies such as mass spectrometry (MS), ion mobility spectrometry (IMS), gas chromatography (GC), optical spectroscopy, have been developed to detect a wide range of explosives and chemical weapons, and, to a lesser extent, biological weapons. The precision of a trace detector is an important factor in minimizing false positive rates. However, collecting enough of a sample to deliver to the trace detector is an important factor in maximizing detection rates, because if a trace sample is not appropriately delivered to the detector, a “non-detect event” may occur. In at least some cases, collection and delivery of samples to these trace detectors is difficult, as methods of collection and delivery must adapt to a wide range of applications and screening scenarios.\nCertain contraband substances—such as explosive or narcotics—are difficult to detect in trace amounts. Security screenings for high traffic areas, such as an airport, attempt to identify and detain individuals involved with contraband substances by collecting samples from the body or clothing. Due to the high traffic nature of these security screenings, the screenings need to be both quick and accurate.\nOne type of sampling method involves using hand-held detection devices. These devices, usually in the form of a “wand”, are easy to handle, provide a suitable distance between the user and the target surface to avoid cross-contamination, and simplify sample transfer to an analytical device. However, these devices generally require the application of a particular amount of force or pressure in order to obtain a sufficient amount of sample from the target surface. Sufficient collection of a sample increases the probability of detecting a trace material within the sample, and the trace detection rate decreases when a user does not apply enough force or pressure against the target surface."}
{"text": "This invention relates to a method of transmitting information between a plurality of stations in a local area network.\nVarious methods are known for sharing network communication channels among stations in a local area network (LAN). One widely used method is known as CSMA/CD, (carrier sense multiple access/collision detect). According to this known method, a station wishing to transmit a message listens until the transmission channel is idle before commencing to transmit an information packet. Furthermore, the station continues to listen to the channel after commencing a transmission and if a collision is detected, that is, more than one station has commenced to transmit an information packet, any station which detects such collision terminates its message transmission and transmits a burst of noise (or garble) such that all other stations can be informed of the collision, terminate message transmission, and wait a random time before attempting to commence a further message transmission.\nFrom Patent Abstracts of Japan volume E-690, page 102, abstract of Japanese Published Application No. 63-187746, there is known a wireless communication system wherein a plurality of stations communicate via a repeater installation. A first frequency is used for transmission from the stations to the repeater installation and a second frequency is used for transmissions from the repeater installation to the stations. If the repeater installation detects a collision on the first frequency transmissions, then it immediately informs all stations of such collision using the second frequency, such that the stations can take appropriate action.\nAlthough the CSMA/CD method has proved satisfactory for transmissions over a cable physical medium, it has been found that a problem arises if it is attempted to apply the method to wireless radio transmissions in that small radio signals may be totally masked by larger signals during reception. This phenomenon is known as the \"capture effect\", and may prevent the detection of collisions in a system such as that disclosed in the aforementioned Japanese patent abstract."}
{"text": "This invention relates generally to gas-insulated power transmission lines, and more particularly to a power transmission line which utilizes two sheaths, one of carbon steel and the other of aluminum for the outer enclosure.\nCompressed gas-insulated power transmission lines are a relatively recent commercial product which is gaining increased acceptance by the electric utility industry as an economical means for transmitting large blocks of electrical energy, and for use in special applications where conventional power transmission lines are inadequate. A typical compressed gas-insulated transmission line is generally comprised of an elongated outer sheath having disposed therein an inner conductor at high potential, with means disposed in the outer sheath to insulatably support the inner conductor within the outer sheath. An insulating gas, typical of which is sulfur hexafluoride at a pressure of 50 psi gauge, is utilized for electrically insulating the inner conductor from the outer sheath. The inner conductor is typically used to transmit energy at voltage ranges from 121-1200 KV. Both the inner conductor and the outer sheath are of good electrically conducting materials such as aluminum.\nOne of the drawbacks of gas-insulated power transmission lines which is hindering even more widespread applications of the lines is associated with the cost of the lines. One of the more costly components of the transmission lines is the outer enclosure utilized to contain the inner conductor and the insulating gas.\nPresent gas-insulated power transmission lines utilize a solid enclosure of aluminum to contain the inner conductor and the insulating gas. Aluminum is preferred because of the low loss it exhibits to the return current flowing through it induced by the conductor currents on the inner conductor. Typically, for grounded aluminum enclosures, the loss ratio (i.e., the ratio of enclosure loss to conductor loss) is of the order of one or less. However, the aluminum material utilized for the enclosure is itself expensive, and large quantities of the material must be utilized.\nOne method attempted in the prior art to minimize the cost of the outer enclosure has been to utilize carbon steel as the outer enclosure. Although carbon steel pipe is generally readily available at less expensive costs, the losses which result with respect to the return current are substantial. The loss ratio of carbon steel is generally of the order of three or higher. Therefore, because of the high-loss ratio, the use of carbon steel for the outer enclosure has generally occurred only when the current rating of a line is very low and the high electrical losses are acceptable."}
{"text": "A dental composition contains a polymerizable monomer, a polymerization initiator, and an inorganic filler, and is most widely used today as a restorative material for repairing fractures of teeth and dental caries. Such a dental composition is required to have the following properties. Specifically, as a cured product obtained after polymerization, the dental composition is required to have sufficient mechanical strength and hardness to serve as a substitute for natural teeth, wear resistance against occlusion of teeth in an oral cavity, surface smoothness and gloss, color matching with natural teeth, transparency, etc. Furthermore, from the viewpoint of handling properties, the dental composition, as a paste which has not yet been polymerized, is desired to have ease of handling for dental clinicians, for example, proper fluidity and high forming property. Particularly, a type of composition to be filled directly in a cavity must have a viscosity low enough to be filled directly from a syringe. If a specific component is added or the blend ratio of the components is adjusted, the resulting composition can produce favorable effects as a dental material.\nThese properties of the dental composition are greatly influenced by the component materials, shape, particle size of an inorganic filler used therein, and further by the content of that filler. For example, when an inorganic filler having a large average particle size of 1 μm or more is used, the filling rate of the filler in the polymerizable monomer can be increased easily and therefore sufficient mechanical strength as a cured product and high handling properties as a paste can be obtained. The use of such an inorganic filler has, however, a drawback in that it is difficult to obtain satisfactory gloss even after final polishing. On the other hand, when an inorganic ultrafine particle filler having an average particle size of less than 1 μm is used, the surface smoothness and gloss after polishing of the cured product and the gloss durability in the oral cavity are improved. The use of such an inorganic ultrafine particle filler has, however, a drawback in that when the inorganic filler is mixed and kneaded with the polymerizable monomer, the viscosity of the resulting paste increases significantly, which makes it difficult to increase the content of the filler. As a result, the mechanical strength of the cured product decreases, and the unpolymerized pasty composition becomes sticky, which reduces the handling properties. When a conventional dental composition, particularly a type of composition to be filled directly in a tooth is used, a reduction in the content of an inorganic filler seems to be a good solution to meet the handling properties requirements. However, when the content of the inorganic filler is reduced, the physical properties such as flexural strength of the resulting composition are low in value. When the content of the inorganic filler is increased to increase the strength, the viscosity of the resulting composition also increases. Such a composition cannot be used for direct filling during dental treatment. Furthermore, a spherical filler can be used to increase the surface smoothness and gloss, but it is difficult to increase the strength while maintaining the handling properties. Under these circumstances, it is difficult to increase the mechanical strength and the surface smoothness and gloss after polishing of the cured product and the handling properties of the paste in a balanced manner.\nPatent Literature 1 discloses a dental restorative material in which an inorganic filler treated with a specific silane coupling agent and a polymerizable monomer having high hydrophobicity are used in combination to achieve high density filling, high strength, high aesthetic quality, and durability. In a dental restorative material disclosed in Patent Literature 2, a mixed filler composed of irregular-shaped inorganic particles, spherical inorganic particles, and inorganic ultrafine particles are used. The irregular-shaped particles having a small average particle size are used in this mixed filler, and acylphosphine oxide is used as a photopolymerization catalyst, so that the resulting dental restorative material has increased surface smoothness and gloss while maintaining high fracture toughness and strength. Patent Literature 3 discloses a dental composition in which a spherical filler having a specific particle size is used to achieve a good balance between the paste properties and the polishability."}
{"text": "1. Field of the Invention\nRecent advances in the general field of molecular biology have made it possible to detect specific genes of clinical and commercial importance. For example, the structures of various genes and gene sequences associated with specific human diseases are known, as are various techniques for detecting the presence of such genes. It is therefore possible to diagnose human disease at the genetic level.\nThe most common technique for detecting a particular gene sequence is hybridization. A particular nucleotide sequence or \"probe\" is marked with a detectable label, typically a radioactive label or chemical modification, and combined with the nucleic acid sample of interest, either in situ as part of intact cells or as isolated DNA or RNA fragments. The sample can be either free in solution or immobilized on a solid substrate. If the probe molecule and nucleic acid sample hybridize by forming a strong non-covalent bond between the two molecules, it can be reasonably assumed that the probe and sample are essentially identical. The probe's detectable label provides a means for determining in a known manner whether hybridization has occurred and for measuring the amount of DNA/RNA sample present. The hybridization technique is of prime importance in basic research directed at understanding the relationship between nucleotide sequences and their function, as well as in diagnostic use to detect known aberrant genes or disease agents such as viruses or bacteria.\nThe main limitation of present gene detecting methods is that they are not sensitive enough and therefore require a relatively large amount of sample to accurately verify the existence of a particular gene sequence. This is not surprising since the detection of a single gene in the entire genetic repertoire of a human being requires locating one part in one to ten million. In fact, most hybridization methods require at least one to ten micrograms of purified DNA, representing a substantial sample of cells, to perform a reliable analysis. This limitation is particularly significant in pre-natal diagnosis of genetic disorders where only a small cell sample can be taken or in identifying infectious agents such as viruses in small tissue samples. Consequently, there is a substantial need for gene detecting methods which will increase the sensitivity of the hybridization assay without sacrificing its specificity.\n2. Description of the Prior Art\nThe hybridization procedure typically includes the initial steps of isolating the DNA sample of interest and purifying it chemically. The DNA sample is then cut into pieces with an appropriate restriction enzyme. The pieces are separated by size through electrophoresis in a gel, usually agarose or acrylamide. The pieces of interest are transferred to an immobilizing membrane in a manner that retains the geometry of the pieces. The membrane is then dried and prehybridized to equilibrate it for later immersion in a hybridization solution.\nA probe labeled with a radioactive isotope is constructed from a nucleotide sequence complementary to the DNA sample by a conventional nick translation reaction, using a DNase and DNA polymerase, although other types of labels can be used. The probe and sample are then combined in a hybridization buffer solution and held at an appropriate temperature until annealing occurs. Thereafter, the membrane is washed free of extraneous materials, leaving the sample and bound probe molecules supported on the membrane. The signal of the bound probe molecules is typically detected and quantified by autoradiography and/or liquid scintillation counting.\nSouthern, J. Mol. Biol. 98:503 (1975) teaches the transfer of DNA fragments to strips of nitrocellulose after the fragments have been resolved by electrophoresis in agarose gels. Immobilization of DNA on diazobenzyloxymethyl cellulose is taught by Noyes and Stark, Cell 5:301 (1975). Alwine et al., PNAS USA 74:5350 (1977) teach a method for detecting specific RNA molecules that are resolved in agarose gels, transferred to diazobenzyloxymethyl paper and then hybridized with DNA probes. Similarly, Reiser et al., Biochem. Biophys. Res. Comm. 85:1104 (1978), teach the detection of small DNA fragments resolved in polyacrylamide gels by immobilizing the same on diazobenzyloxymethyl paper and then hybridizing them to DNA probes.\nBoth Bittner et al., Anal. Biochem. 102:459 (1980) and Stellwag and Dahlberg, Nucl. Acid Res. 8:299 (1980) teach the use of an electrical current to transfer DNA and RNA fragments from either agarose or acrylamide gels to diazobenzyloxymethyl paper. Similarly, electrophoretic transfer of nucleic acids to diazophenylthioether paper is taught by Reiser and Wardale, Eur. J. Biochem. 114:569 (1981).\nA New England Nuclear (Boston, Mass.) technical bulletin entitled \"Gene Screen Hybridization Transfer Membrane Instruction Manual\", by D. J. Green and D. R. Rittenbach (1982), teaches a modification of two different techniques (\"electrophoretic transfer\" and \"capillary wicking\") for immobilizing DNA and RNA samples on a porous substrate, specifically nylon-base membranes.\nNucleic acid probes have heretofore been labeled separately with either tritium or radioactive phosphorus [.sup.32 P] by nick translation, Rigby et al., J. Mol. Biol. 113:237 (1977) or with biotinylated uridine, Narayanswami, Hutchison and Ward, J. Cell Biol. 95:74a (1982). Other references, besides Rigby et al., identifying radioactive labels for DNA probes include Wahl et al., U.S. Pat. No. 4,302,204 (.sup.32 P); Falkow et al., U.S. Pat. No. 4,358,535 (.sup.32 P, .sup.3 H, .sup.14 C); and Axel, et al., U.S. Pat. No. 4,399,216 (.sup.32 P) A New England Nuclear technical bulletin, Vincent et al., \"Preparation of DNA Labeled With High Specific Activity [.sup.35 S]-Deoxyadenosine 5' - [.alpha.-Thio] Triphosphate; the Use of .sup.35 S-Labeled Nucleic Acids as Molecular Hybridization Probes\" (1982), teaches that a hybridization probe labeled with [.sup.35 S]-deoxyadenosine 5' [.alpha.-thio] triphosphate is qualitatively indistinguishable from a conventional .sup.32 P labeled probe.\nWahl et al., U.S. Pat. No. 4,302,204 teach the use of ionic polymers, particularly dextran sulfate, to increase the local concentration of nucleic acids in hybridization reactions and in turn increase the signal of the DNA sample of interest. It also describes the use of a depurination step to increase the efficiency with which very large nucleic acid segments are transferred from a gel to a solid membrane.\nA method of incorporating a phosphorothioate analog of deoxynucleotides, namely .sup.32 S, into DNA polymers, using known DNA polymerases (E. coli DNA polymerase I and E. coli DNA polymerase III) is described by Kunkel et al., PNAS USA 78:6734 (1981). This reference mentions the use of thionucleotides to induce site-specific mutations as an application of the method. Putney et al., PNAS USA 78:7350-7354 (1981) teach that a protective effect against enzymatic degradation occurs following the incorporation of a thionucleotide (.sup.32 S) into DNA polymers.\nThe use of a polyfunctional disulfide compound to cross-link protein molecules is taught by Kotani et al., U.S. Pat. No. 4,287,345.\nThe present invention improves prior methods of detecting DNA genes and gene sequences by amplifying the detectable signal generated by bound probe molecules, thereby providing an accurate and reliable method for detecting and quantifying particular nucleotide sequences, even in relatively small samples. Accordingly, the present invention includes as its objects the following:\n(1) to increase the mass of complementary DNA bound to specific gene sequences immobilized on a membrane; PA0 (2) to increase the signal-to-noise ratio produced in hybridization reactions involving immobilized DNA; PA0 (3) to improve the resolution of bands visualized by autoradiography of radioisotopically-labeled DNA; PA0 (4) to produce radioisotopic probes which will provide the same intensity of autoradiographic signal as obtained with .sup.32 P-labeled probes, but with less danger of radiation exposure to the worker due to the decreased emission energy of the probe; PA0 (5) to improve existing protocols for hybridizations, especially filter hybridizations, with .sup.35 S-labeled probes and reduce both the direct cost of the reaction mixture and the labor cost of the procedure; and PA0 (6) to produce a system of signal amplification which is not dependent upon the use of radioisotopes and which is compatible with a variety of probe-labeling systems.\nOther objects will be apparent from the Description of a Specific Embodiments."}
{"text": "The present invention relates generally to intake structures for internal combustion engines, and more particularly, to a vane which is incorporated in the intake passage of internal combustion engines for producing advantageous swirl in the combustion chambers of the engine during operation.\nIn the field of internal combustion engines, it is known that, if the intake gas induced into a combustion chamber is provided with a swirling motion, the resulting turbulence assists the propagation of flame in the combustion chamber increasing the speed of flame propagation, thereby improving the performance of the engine. In addition, the anti-knock and driving characteristics of the engine, and the quality of the exhaust emissions thereof, are improved.\nConventionally, the inlet port of an internal combustion engine, especially a gasoline engine, is formed so as to minimize the resistance to the flow of intake gases and thereby increase the volumetric efficiency and maximum output power of the engine at high load operation. When such an internal combustion engine is operating at high load, although the flow or suction resistance of the inlet port is low, sufficient turbulence and swirl is caused in the combustion chamber for the flame propagation speed to be kept at an acceptable level.\nHowever, during low and medium load operation, in which the amount of inlet gases is relatively low, a good level of swirl and turbulence of the inlet gases does not occur, due to the low flow resistance of the inlet port, and because the velocity of the inlet gas entering the combustion chamber from the inlet port is low. Thus combustion speed is sometimes not high enough, and, especially if the engine is operating in a lean mixture condition, or is provided with a substantial quantity of exhaust gas recirculation, good operating performance of the engine may be lost.\nConventionally, it has been realized that the generation of swirl in the cylinder is strongly affected by the direction of entry of inlet gases thereinto, and their degree of swirling as they enter. Therefore, in the past, various types of inlet ports have been proposed, with a view to enhancing swirl. One of these has been the so-called helical port, which is provided with a helical surface leading to the valve seating surface and extending helically around the axis of the valve seating surface. This type of port imparts a swirling action to the inlet gases before they pass the inlet valve.\nHowever, there is a significant disadvantage in the performance of such a helical inlet port, which limits the usefulness thereof. If the helical surface is provided as a long and deep helix, so that it imparts a very strong swirling action to the inlet gases, then, at high load operation, when a large volume of inlet gases is required to be passed by the inlet port, the volumetric efficiency of the engine will be degraded, and thus its maximum output power will suffer, due to the increase of the resistance of the inlet passage to the flow of inlet gases. On the other hand, if the helical surface is provided as a short and shallow helix, so that it imparts a rather weak swirling action to the inlet gases, it will not function well at low load operation of the engine. Accordingly, such a helical port has not fully met the need for generating swirl in intake gases.\nThe present invention results from the realization that, if the flow of inlet gases is specifically directed to the start of the helical surface of a helical inlet port during low and medium load operation of the engine, without greatly increasing flow resistance of the inlet port to inlet gases, swirl and turbulence in the combustion chamber would be advantageously improved, and combustion speed would be increased, without sacrificing engine output power.\nIt is, therefore, an object of the present invention to provide an apparatus in the inlet passage of an internal combustion engine by which the mixture gas is directed to the start of the helical surface of a helical inlet port, particularly during low and medium load operation of the engine, without greatly increasing flow resistance of the inlet passage."}
{"text": "1. Field of the Invention\nThe presently disclosed subject matter relates to a technique that retrieves and presents a case image similar to a diagnostic object image (image of object to be diagnosed) on the basis of an amount of feature (feature amount) extracted from the diagnostic object image.\n2. Description of the Related Art\nIn Japanese Patent Application Laid-Open No. 2004-5364, an ROI (region of interest) is specified on image data; information such as image data which has a similar feature and a diagnostic result of the similar image data are retrieved from a database and displayed.\nIn Japanese Patent Application Laid-Open No. 2001-117936, a region of interest is extracted from a diagnostic image, and a similarity search is performed using the amount of feature of the region of interest.\nJapanese Patent Application Laid-Open No. 2008-257292 shows another example of similar image search system.\nJapanese Patent Application Laid-Open No. 2008-245719 shows an example of a technique for extracting a lesion region on the basis of position information of a region of interest.\nT. F. Cootes, G. J. Edwards, and C. J. Taylor “Active Appearance Models”, In Proc. 5th European Conference on Computer Vision, Springer Germany, volume II: pp. 484-498 (1998), shows a method of acquiring shape information and texture information of a lesion part in the ROI (region of interest)."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory device, particularly a dynamic semiconductor memory device (DRAM).\n2. Related Background Art\nIn a related DRAM, a memory cell is composed of an MOS transistor and a capacitor. The scale-down of the DRAM has been remarkably advanced by the adoption of a trench capacitor structure and a stacked capacitor structure. At present, the cell size of a unit memory cell is scaled down to an area of 2 F×4 F=8 F2, where F is a minimum feature size. Namely, the minimum feature size F decreases with the advance of generation, and when the cell size is generally taken to be αF2, a coefficient α also decreases with the advance of generation. Thus, at the present of F=0.18 μm, α=8 is realized.\nIn order to hereafter secure the trend of cell size or chip size which is the same as before, it is demanded to satisfy α<8 in F<0.18 μm and further satisfy α<6 in F<0.13 μm, and together with microfabrication, the formation of cell size of the possible small area becomes a large problem. Accordingly, various proposals for decreasing the size of the one memory cell with the one transistor and one capacitor to 6 F2 or 4 F2 are made. However, practical use is not easy since there are a technical difficulty that the transistor has to be a vertical type, a problem that electric interference between adjacent memory cells increases, and in addition difficulties in terms of manufacturing technology including fabrication, film formation, and the like.\nOn the other hand, some proposals for a DRAM in which a memory cell is composed of one transistor without using a capacitor are made as mentioned below.\n(1) JOHN E. LEISS et al, “dRAM Design Using the Taper-Isolated Dynamic Cell” (IEEE TRANSACTIONS ON ELECTRON DEVICES, VOL. ED-29, NO. 4, APRIL 1982, pp 707-714)\n(2) Japanese Patent Laid-open Publication No. H3-171768\n(3) Marnix R. Tack et al, “The Multistable Charge-Controlled Memory Effect in SOI MOS Transistors at Low Temperatures” (IEEE TRANSACTIONS ON ELECTRON DEVICES, VOL, 37, MAY, 1990, pp 1373-1382)\n(4) Hsing-jen Wann et al, “A Capacitorless DRAM Cell on SOI Substrate” (IEDM93, pp 635-638)\nA memory cell in (1) is composed of MOS transistors, each of which has a buried channel structure. Charge and discharge to/from a surface inversion layer is performed using a parasitic transistor formed at a taper portion of an element isolation insulating film to perform binary storage.\nA memory cell in (2) uses MOS transistors which are well-isolated from each other and uses a threshold voltage of the MOS transistor fixed by a well potential as binary data.\nA memory cell in (3) is composed of MOS transistors on an SOI substrate. A large negative voltage is applied from the SOI substrate side, and by utilizing accumulation of holes in an oxide film of a silicon layer and an interface, binary storage is performed by emitting and injecting these holes.\nA memory cell in (4) is composed of MOS transistors on an SOI substrate. The MOS transistor is one in terms of structure, but here a structure, in which a reverse conduction-type layer is formed on top of the surface of a drain diffusion region, whereby a P-MOS transistor for write and an N-MOS transistor for read are substantially combined integrally, is adopted. With a substrate region of the N-MOS transistor as a floating node, binary data are stored by its potential.\nHowever, in (1), the structure is complicated and the parasitic transistor is used, whereby there is a disadvantage in the controllability of its characteristic. In (2), the structure is simple, but it is necessary to control potential by connecting both a drain and a source of the transistor to a signal line. Moreover, the cell size is large and rewrite bit by bit is impossible because of the well isolation. In (3), a potential control from the SOI substrate side is needed, and hence the rewrite bit by bit is impossible, whereby there is a difficulty in controllability. In (4), a special transistor structure is needed, and the memory cell requires a word line, a write bit line, a read bit line, and a purge line, whereby the number of signal lines increases."}
{"text": "Nowadays, most video signals have been recorded on tapes in an analog form. The quick uptake of digital video for playback, for example using DVD, has as a consequence that demand for recording video material on a digital recording medium also increases, for example using a recordable optical disc, like DVD+RW, or hard disc based recording devices. A commonly used digital encoding format is MPEG, in particular MPEG2. For recording an analog video signal, the user uses an analog reproducing device, such as a VCR and Camcorder to generate the analog video signal and uses an (MPEG) encoder to encode this signal to a digital format for storing on a digital storage medium. The MPEG encoder may be a stand-alone device, or integrated in another device, like the VCR, television, etc. Via a digital video communication system, such as the P1394 bus, the encoded signal can be supplied to a decoder/renderer or to a digital storage device for subsequent supply to a decoder/renderer. To enable decoding/rendering the encoded signal must comply with the encoding/transmission standard, such as MPEG. MPEG video encoding/transmission is, among others, described in International Standard ISO/IEC 13818-2: 1995(E) Information Technology; Generic coding of moving pictures and associated information: Video and International Standard ISO/IEC 13818-1: 1996(E) Information Technology; Generic coding of moving pictures and associated information: Systems. Especially timing of the stream of digital video is important. To this end, digital video encoding/transmission standards specify clock signals that enable clocks in the encoder and decoder to run in a synchronized manner. These clock signals will be referred to as program clock reference stamps (PCR). Additionally, timing signals are supplied for the individual frames to enable a rendering device to render the frames at the correct time. These timing signals will be referred to as presentation time stamps (PTS). Moreover timing signals can be present that indicate to the decoder when the MPEG decoding can start. These signals will be referred to as decoding time stamps (DTS). The analog or digital video signal supplied to an encoder in addition to the video signal also includes a frame signal (vsync) that indicates frame boundaries in the video signal. The frame signal is thus synchronous with the video frames. The PTS signals in the encoded stream are usually inserted in the encoded video stream synchronous with the arrival of the vsync signal.\nFor MPEG2 transmission as specified for DVB (Digital Video Broadcast), the timing for the PCR and PTS are based on a same clock with a defined frequency and tolerance. For example, the clock defined by MPEG runs at 27 MHz. and should deviate less than 30 ppm (parts per million). Usually, the frame rate of a video signal is 25 Hz. or 30 Hz. and the vsync pulse has the same frequency. A straightforward conversion of an analog video signal having a frame rate of exactly 25 Hz. will result in an MPEG stream with PTS values present once every 40 ms. However, the deviation of the frame rate of a video signal reproduced with analog reproducing devices may be up to 100 ppm. To be able to obtain a digital video signal which has a PTS once every frame period, the MPEG encoder could be locked to the vsync of the video input signal. The encoded digital video signal thus obtained will have a PCR which deviates 100 ppm and is therefore not MPEG compliant. Such an encoded video stream should not be supplied through a network to a decoder/renderer as the decoding device may not be able to decode the signal. On the other hand, directly using the MPEG clock and ignoring too large deviations in frame rate may result in rendering problems, such as having to skip or repeat frames."}
{"text": "1. Field of the Invention\nThe invention relates to a method and a corresponding device for vehicle communication between a first vehicle and a second vehicle.\n2. Related Art\nIn road traffic, vehicle drivers often communicate using the vehicle horn or light signals. However, these signals are often difficult to interpret for the vehicle driver of the vehicle to whom the signal is directed.\nEP 1 515 292 A1 discloses a communication system for a vehicle, in which light is emitted by a headlight or a rear light of the vehicle. A signal from a modulator is superimposed on the light. The modulated light is received by a light receiving unit of a second vehicle. The received modulated light is demodulated by a decoder for display on a display panel.\nDE 10 2010 007 259 A1 discloses a method for transmitting information relating to a defect in at least one part of a motor vehicle. In the method, a defect in at least one part of a second motor vehicle is first of all determined by at least one occupant of a first motor vehicle. A first item of information, which uniquely identifies the second motor vehicle is input and a second item of information, which identifies the defect, is input. A first message is transmitted to a central distribution unit, the first message comprising the first item of information and the second item of information. The first item of information and the second item of information are also processed in the central distribution unit and a second message is transmitted from the central distribution unit to at least one contact address."}
{"text": "Over the years, one of the most successful floor covering products has been vinyl tile flooring. The tile is long wearing and attractive, particularly when clear vinyl is laminated over wood or other decorative and eye pleasing scenes and materials. A laminate containing clear vinyl over a design, with additional vinyl and an adhesive system, allows for the mounting of vinyl tile on subfloors. These tiles are ideally suited for high traffic areas. It is essential to present a very favorable impression to consumers as they enter a store, beginning with clean and polished floors. Vinyl tile is presently the best floor covering assembly when factors such as cost, permanence, cleanliness, appearance, and durability are all considered. For that reason, vinyl tile is the preferred floor for a wide variety of stores, restaurants and other establishments.\nTypical vinyl tile which is commercially available is extremely difficult to remove once it has been installed. Typically, removal requires workers to dig up the floor using a chisel to remove the tile, and then apply a patching material Because the patching must dry, a 24 hour waiting period is necessary. Then the crew returns and the subflooring can be prepared and new tile installed. This is an extremely expensive process, is very unsightly, and calls for significant down time for the particular area of the flooring. Up until the present time, it has not been possible to change the design of a floor covered with vinyl tile without this expense and inconvenience.\nOne alternative which has been proposed is to provide exterior designs which are placed on top of already existing vinyl tile. This alternative presents several significant difficulties. First, if the design is to be permanent and have the strength, resilience and durability of vinyl tile, it must have the same permanence of mounting. Thus, when the message or design is to be changed, the same difficulties in removing an overcoating of some form is realized as would be with the tile originally on the floor area. If the surface pattern is not as durable as vinyl tile, then ordinary pedestrian traffic will disturb or destroy the coating, and at the least make it quite unsightly in a relatively short period of time. Thus, continued replacement and repair increase costs prohibitively. Finally, there is a danger that any floor covering which is significantly different in height, for example, than the adjacent area will present special risks and hazards to unsuspecting pedestrians.\nAccordingly, it is an object of this invention to provide a floor covering assembly which permits a defined area to be conveniently removed and replaced with alternative floor covering.\nAnother object of the present invention is to provide a safe, stable, highly resistant floor covering which can be removed at will and replaced with similar flooring and yet which can be left on the floor for extended periods of time without concern for durability, eye appeal or other factors.\nOther objects will appear hereinafter."}
{"text": "1. Field of the Invention.\nThe present invention relates to a rip fence for use with a table saw and, in particular, to a rip fence that has a dual locking mechanism for securing the fence in position.\n2. Scope of the Prior Art.\nTable saws of various designs include a frame and an upper surface having an opening through which a blade extends. The upper surface supports a material, such as wood, as it is pushed towards the blade for cutting. Table saws have traditionally been stand-alone pieces of equipment that are used in workshops. Because most stand-alone table saws are used in large open spaces, the dimensions of the table saw can be large enough to support materials of various sizes for cutting. Tables that are used in workshops are not readily moved from one location to another. At construction sites, it is useful to have a table saw that can be easily moved from one location to another or between different construction sites. Accordingly, smaller and portable table saws have been designed. These table saws use lighter materials and have reduced the dimensions of the frame and the upper surface to reduce the size and weight of the table saw thereby creating the desired portable table saw. The portable table saws also use a separable folding table to hold the table saw at the correct height for operation.\nAn integral part of most table saws is a movable rip fence. The rip fence is positioned perpendicular to the upper cutting surface and is aligned generally parallel to the position of the blade. Most rip fences are designed to span between the front edge and the rear edge of the upper cutting surface. Moreover, the rip fence can be positioned in any position along the upper surface and is secured in such a position by a locking mechanism. The rip fence is usually secured in a position a given distance away from the blade and can be put on either side of the blade. When the rip fence is secured in a position, the wood can be pushed against the side of the rip fence to slide on the upper surface to be cut by the blade. Thus, accurate and straight cuts can be made.\nRip fences of various different designs are readily known in the prior art. In addition, different methods of securing the rip fence at a given location on the upper surface, including the use of a locking mechanism, are known. One of the primary objectives for the rip fence is for it to be generally parallel to the blade so that the most accurate cuts can be made. It is known that an effective way to make the rip fence parallel with the blade is to force it to be perpendicular to the front and rear edge surfaces of the table saw. Thus, many rip fences use a front clamping mechanism and a rear clamping mechanism. The front clamping mechanism presses against the front edge surface and the rear clamping mechanism presses against the rear edge surface. In this way, the clamping mechanism aligns the rip fence perpendicular to the edge surfaces and parallel to the blade.\nIt is also helpful to properly align the rip fence if the front and rear clamping mechanism move simultaneously. To achieve this, rip fences typically include a handle which is connected to the front clamping mechanism and a rod that connects between the handle and the rear clamping mechanism. When the handle is in a first and released position, the front locking mechanism is positioned in a released position relative the front edge surface and the rod reduces pressure and allows the rear locking position to move into a released position relative the rear edge surface. When the handle is in a second and locked position, the front locking mechanism is positioned in a locked position along the front edge surface and the rod pushes the rear locking mechanism into the locked position. As the handle moves from the released to the locked position, the front and rear clamping mechanism engage with the front and rear edges at the same time.\nDifferent types of front and rear locking mechanisms are known. Some front locking mechanisms include a cam that is connected to a handle such that when the handle is put into the locked position, the cam is pushed against a surface on the front edge of the table saw to secure the rip fence in position. U.S. Pat. No. 4,846,036 to Metzger, Jr., et al., reveals a handle that moves a spring-biased clamping portion. In the locked position, the clamping portion is biased against a vertical wall in a slot that is proximate the front edge surface of the upper surface. In the released position, the handle overcomes the biased clamping portion so that it is released from the wall thereby allowing the rip fence to move along the upper surface. As the rip fence moves along the upper surface, the clamping portion moves through the slot.\nOther rip fences are provided with different mechanisms to move smoothly along the upper surface. As described, slots can be provided along the upper surface of the table saw through which a portion of the rip fence, like a clamping portion, can slide. In addition, tubular rails can be provided on the front and rear sides of the table saw. In those embodiments, the rip fence include a portion that surrounds the tubular rail so that the rip fence can slide into a selected position. Other types of rails can be provided underneath the upper surface. In addition, the rip fence can include hook portions, tabs or other mechanisms that move relative a part of the upper surface or the frame so that the handle can lock the rip fence in a location.\nBecause of the rip fences\"\" designs, the handles used to lock the rip fence create inefficiencies. Many such handles are quite long so that their rotation provides enough torque to adequately secure the rip fence in a selected position. Those handles often protrude significantly out from the front edge of the saw where the handle can get in the way of efficiently using the table saw. The large handles can also catch on clothing and other items. In order to reduce the size, other rip fences reorient the handle\"\"s location on the rip fence. Other prior art methods alter the direction that the handle rotates relative the table saw. Those alternative embodiments are not as easy to use.\nIn view of the foregoing, it is an object of the present invention to develop a portable saw that has a rip fence to overcome the deficiencies of the prior art. What is needed is a rip fence that moves between multiple positions on the table saw, that can be secured easily in any such position, and that does not unnecessarily increase the weight or general dimensions of the table saw. It is also desired to have a rip fence that is aligned parallel to the blade when it is locked into position. Moreover, the rip fence should reduce the dimensions of the handle so that it does not protrude significantly in front of the table saw.\nThe present invention relates to a table saw that includes a rip fence that is slidable relative to the table saw. The table saw used with the rip fence includes a frame portion and upper cutting surface. A blade extends through an opening in the upper surface to cut wood or other material. The upper surface supports the wood as it is pushed on the table saw and cut by the blade. Front and rear edge surfaces extend perpendicularly from the front and rear edges of the upper surface.\nThe rip fence of the present invention can be used with a table saw that has movable auxiliary table. The auxiliary table can move between a first position adjacent the upper surface and a variable second position separated from the upper surface according to any known methods. One such method includes rails that are secured to the front and rear edge surfaces of the auxiliary table and that slide along brackets on the rear and front edge surfaces of the upper surface.\nThe rip fence moves between positions relative the upper surface and can be locked in any such position such that the rip fence is parallel to the blade. This arrangement ensures that accurate cuts are made. To move relative the upper surface, the rip fence engages with a slot which is parallel to the front edge surface formed on the upper cutting surface of the table saw or the upper surface of the front rail. The rip fence can also engage with an groove formed in the rear edge surface of the upper cutting surface or the outer edge of the rear rail.\nThe rip fence extends between the front edge surface and the back edge surface of the upper cutting surface and includes a main casing that is slightly shorter than the span of the entire rip fence. A front clamping mechanism is provided that is connected to the front end of the casing by a frame. A movable handle is pivotally connected at one end of the handle\"\"s upper surface to the frame and rotates between a lower locked position and an upper released position. When the handle is in the locked position, the rip fence is held securely in position relative the upper cutting surface so that the rip fence is generally parallel to the blade. When the handle is in the released position, the rip fence can slide relative the upper cutting surface and longitudinally through the slot and groove.\nThe front clamping mechanism includes a rotatable cam that has a first end and a second end. In addition, a movable link having an upper end and a lower end is provided. The link is pivotally connected at its upper end to a distal end of the handle\"\"s upper surface and pivotally connected at the lower end to the first end of the cam. The cam\"\"s second end is pivotally connected to the frame. A vertically extending activating plate is partially contained within the frame such that its upper end is in contact with the cam\"\"s second end. The lower end of the activation plate extends below the lower edge of the front clamping mechanism. The activation plate can be pivotally connected to the handle frame at a point below the upper end.\nThe front clamping mechanism also includes a horizontally extending clamping plate that engages with the slot. The clamping plate includes a passage proximate a first end and a flange, or lip, at an opposing second end. The lower end of the activation plate extends through the passage. The lip engages with the slot so that as the rip fence moves along the upper cutting surface the lip slides longitudinally through the passage. The lip moves laterally within the slot between a locked position and a released position. In the locked position, the lip presses against the walls of the slot so that the rip fence does not slide. In the released position, the lip is free of the slot\"\"s walls so that the rip fence can slide to a desired position. The clamping plate moves between the locked position and the released position as the handle moves between its locked and released position. As the handle moves between positions, the link rotates the cam and therefore the activation plate. As the activation plate pivots, its lower end laterally moves the clamping plate between the locked and released positions.\nThe rip fence of the present invention also includes a rear clamping mechanism that has a second activation plate. An upper end of the rear activation plate is pivotally connected to a rear frame. A portion, or tab, is provided at a lower end to engage with the groove provided along the rear edge surface of the table saw. A spring-loaded rod maintains contact between the front activation plate and the upper end of the rear activation plate. The second activation plate, and therefore the portion, move between a locked position and the released position as the handle moves between those positions. As the handle moves between those positions, the first activation plate rotates about its pivot point thereby moving the rod. As the rod moves, it forces the rear activation to rotate about its pivot point so that the tab moves between the locked position and the released position. In the locked position, the edges of the portion engage with the walls of the groove so that the rip fence is secured in a position and is parallel to the blade. In the released position, the portion is removed from the groove so that the rip fence can move along the upper surface.\nThe front clamping mechanism, or the casing, can include gliding blocks that slide through the slot. The glide blocks are arranged on either side of the clamping plate and permit the rip fence to move smoothly through the slot between various positions. In addition, the rip fence can include a roller assembly that allows the rip fence\"\"s rear clamping mechanism to move smoothly.\nIn view of the foregoing, a rip fence is provided that moves smoothly between various positions along the upper cutting surface. The multiple pivot points in the front clamping mechanism also provide a handle that is easy to use and that does not protrude out from the table saw. The clamping plate and rear activation plate also provide effective mechanisms to secure the rip fence in a position so that the rip fence does not move during operations of the table saw. In addition, the front and rear clamping mechanisms align the rip fence in a generally parallel arrangement with the blade when they are moved into the locked position.\nThese and numerous other features and advantages of the present invention will become readily apparent from the following description, the accompanying drawings and the appended claims."}
{"text": "The present invention relates to methods and apparatus for identifying features in an image. More particularly, but not exclusively, the present invention relates to methods and apparatus for fitting models to an image to identify features in the image.\nAutomated image analysis may be used to reduce manual workloads in determining properties of images. Image analysis is used in a variety of different fields to identify features of images. One area in which image analysis techniques are used is to process medical images where the identification of anatomical features in an image can provide useful information to medical practitioners. Often many images are required to be processed, for example in a screening program, and each image should be processed as accurately as possible.\nImages of a particular anatomical feature obtained from different patients vary considerably. This may particularly be the case when an anatomical feature is in a diseased state. It can therefore be difficult to accurately identify anatomical features in images given such variance. An existing approach to identifying anatomical features in an image is based upon statistical modelling of an anatomical region of interest. Statistical models are generated from a training set of images in which features within the anatomical region of interest have been manually identified and a statistical model indicates a mean model and an acceptable degree of variation from the mean model, as indicated by the training set. Where a small number of training images are available, a statistical model can often fail to identify particular anatomical features in previously unseen images due to the failure of the statistical model to capture a true indication of acceptable variation.\nAdditionally, existing approaches for fitting statistical models to previously unseen images can be inaccurate, particularly when an anatomical feature is in a diseased state in the previously unseen image. A robust and accurate method of modelling structures in images and fitting such models to previously unseen images is therefore desirable.\nIt is an object of some embodiments of the present invention to obviate or mitigate at least some of the problems set out above."}
{"text": "1. Technical Field\nThe present disclosure generally relates to a touch device, and more particularly to a touch device and a driving method of a touch panel thereof.\n2. Related Art\nWith developing of technology, most electronic devices, such notebooks, cell phones, portable multimedia players, and so on, usually have touch panels as the new generation input interfaces to replace the traditional keyboards. The touch panels can be categorized into the resistive touch panels, the capacitive touch panels, the infrared light touch panels, and the supersonic touch panels, wherein resistive touch panels and the capacitive touch panels are most common products. In generally, when the touch panels operate, the touch panels scan the whole touch area thereof to sense the touch points.\nFor the projected capacitive touch panel, the larger the size of the panel is, the more the number of the sensors is, and the longer the time which the touch panel scans once is. In some condition, the user may not touch the whole touch area of the touch panel when the touch panel is used. Furthermore, assuming the touch panel is applied in the people interacting education system or the public system under the condition that the whole touch area of the touch panel is available, when user operates the touch panel, the user may be interfered by the operation of the other user, or may interfere the operation of the other user. To maintain the whole touch area of the touch area available not only decreases the efficiency, but also introduces interference."}
{"text": "We humans are social creatures, and have socialized, worked, and otherwise collaborated in both small and large groups for thousands of years. Many of these groups or teams have collaborated to achieve great feats for the good of all mankind, feats including building the great pyramids, circumnavigating the globe, establishing new nations, building the world's first flying machine, winning world wars, producing new medical treatments, sending a man to the moon, to name a few. Pick any great achievement and the odds are that there is a great team behind it. All this suggests that the future of the world depends on the work of great teams.\nThe present inventors have recognized that there are at least three problems that are inhibiting formation of more great teams.\nThe first problem is that many teams fail or underperform because of poor interactions between group members and/or because one or more group members lacked essential attributes needed for the group to succeed at whatever mission it takes on. The problem is rampant and generally stems from the fact that teams are built largely using human mental processes. Which are heavily influenced or completely dictated by traditional organizational or social hierarchies; inherent or implicit social, racial, gender, ethnic, or other conscious and unconscious biases within team organizers; and an inability of team organizers to objectively process real data. Wherever these failures occur they are not only costly in terms of lost time and money of businesses and other organizations relying on team performance, but also costly in terms of the wasted effort and emotional stress on team members who joined the teams to achieve something bigger than they could do alone.\nThe second team formation problem recognized by the present inventors is in the context of education, for example within universities, colleges, high schools, and so forth. In this context, many teams are formed for group projects and study groups, etc. The groups are typically small in number and formed during class in rapid ad hoc ways with consideration of only a few of the many relevant variables. For example, teams for group projects are often based largely on seating proximity or alphabetical order, resulting in widely variant team efficacy. Moreover, students enrolled in any given class also constitute a team. These classroom teams are typically self-organized by students or school administrators based on teacher or scheduling preferences, ignoring numerous individual student traits that affect the entire group learning experience within every class.\nThe third problem recognized by the present inventors is related to formation of teams or groups in the context of temporary events, such as conferences and conventions, where hundreds or thousands of people may come together to learn or confer around a particular topic, such as education, healthcare, or technology. Many of the participants come not only to learn from speakers at these events, but also to meet new people and form new relationships that can lead to professional and personal collaborations. In most of these conferences and conventions, the peer-to-peer interactions are largely serendipitous, that is, uncoordinated in any systematic or intelligent way around particular user interests or desires, leaving participants to rely on ad hoc introductions and happenstance or to generally fend for themselves in coordinating advantageous meetings. As a result, many of the participants miss out on opportunities to further leverage their investment in money, time, and energy to attend conferences and conventions.\nBeyond these problems in organizational, educational, and event team or group formation contexts, the present inventors also recognize that, in broader terms, how well we get along as humans at every level—within friendships, partnerships, and local, national, and international communities—depends on our collective capacity to collaborate and work effectively. While the internet has powered unprecedented interconnectivity, global communications, and informational sharing, it has also exponentially amplified our capacity to behave according to implicit and explicit human biases at every level of interaction, giving rise to insular social, economic, and political echo chambers that regard each other as enemies to oppose rather than friends to commune with. This insularity breeds increasing rigidity in thought and behavior that threaten our capacities not only to understand each other and live peacefully together, but also to collaborate with sufficient diversity to address our grand global challenges.\nAccordingly, the present inventors have recognized a need for better ways of bringing people together into small groups or teams to achieve desired objectives."}
{"text": "1. Statement of the Technical Field\nThe invention is directed to integrated circuit (IC) devices, and more particularly IC devices including inductor structures formed therein.\n2. Description of the Related Art\nReactive components, such as capacitors and inductor, are important components for implementing various types of signal processing devices, including low-noise amplifiers (LNAs), voltage controlled oscillators (VCOs), filters, and impedance matching networks, to name a few. As a result, a significant number of components are typically needed to implement such devices. However, the number of components for such devices can be reduced by using an IC including inductors and capacitors formed therein. As a result, signal processing devices can be formed in a compact package, reducing overall device size.\nThe performance of inductors, including inductors formed within an IC (monolithic inductor), is typically evaluated based on the quality (Q) factor. An ideal inductor will generally be lossless irrespective of the amount of current through the winding. However, inductors, including monolithic inductors, typically have winding resistance from the electrically conductive materials used for forming the coils. Since the winding resistance appears as a resistance in series with the inductor, it is often called the series resistance. The inductor's series resistance converts electrical current through the coils into heat, thus causing a loss of inductive quality. The Q factor measures this loss (i.e., the inductor's efficiency), as Q factor for an inductor is the ratio of its inductive reactance to its resistance at a given frequency. In general, the higher the Q factor of the inductor, the closer it approaches the behavior of an ideal, lossless, inductor. In general, the Q factor of an inductor can be computed from Q=(ωL)/R, where R is its internal electrical resistance of the inductor and ωL is capacitive or inductive reactance at resonance.\nIn the case of monolithic inductors, the Q factor is primarily limited by conductor losses arising from conductor resistances, the conductive silicon substrate, and parasitic substrate capacitances (which lower the inductor self-resonant frequency). Additionally, time-varying magnetic fields can penetrate the silicon substrate and cause eddy currents as per Lenz's law, thus resulting in power loss. Furthermore, eddy currents create their own magnetic fields that oppose those of the monolithic inductor and effectively decrease the inductance of the inductor."}
{"text": "There is conventionally known an image forming apparatus provided with a main casing, an intermediate transfer belt, a plurality of process units disposed below the intermediate transfer belt, and a drawer supporting the plurality of process units and drawable from the main casing. In this configuration, the process units each integrally include a photosensitive drum and a developing unit having a developing roller for supplying toner to the photosensitive drum and a toner chamber, and can each be replaced with new one by drawing the drawer from the main casing."}
{"text": "This invention relates generally to switch constructions, and deals more particularly with a combination kill switch and indicator light for an internal combustion engine of the type driven by a magneto ignition system.\nApplicant is aware of the following prior art patents as bearing some relevancy to the disclosure provided in the subject application. U.S. Pat. No. 4,347,417 illustrates a rocker switch having an internal lamp with a coil compression spring provided for continuous connection with that lamp.\nAnother prior art patent of relevance to the subject disclosure is U.S. Pat. No. 4,454,400 wherein a lamp is provided in a rocker of a switch that also includes means for providing a resistor in that rocker as well. Both these U.S. Pat. Nos. 4,347,417 and 4,454,400 utilize a spring element provided on a depending post on the rocker so as to afford electrical connection between the lamp, the resistor and the rocker itself, and a movable contact that is engaged by the spring so as to urge that movable contact downwardly against one or more fixed contacts in the bottom wall of the switch case.\nThe general purpose of the present invention is to provide a unique application for a switch of the type disclosed in the above-identified prior art patents, and more particularly to provide a unique switch construction that is well adapted to serve as both a low oil indicator light and a kill switch for an internal combustion engine of the type having a magneto ignition system."}
{"text": "The processing of many food and beverage products requires subjecting the product to high temperatures for a period of time, thereby sterilizing the product. Traditional techniques, such as canning, expose the product to high temperatures for extended periods of time. Canning not only has a high energy cost, but also subjects the product to thermal stress, which frequently has a negative impact on the flavor and texture of the product.\nMore recent techniques involve aseptic processing, in which liquid foods and beverages are sterilized outside the package using an ultra-high temperature process that rapidly heats, and then cools, the product before packaging. The processing equipment allows the time (generally 3 to 30 seconds) and temperature (195° to 300° F.) to be tailored to place the least amount of thermal stress on the product, while ensuring safety. The rapid heating and cooling aspects of the aseptic process substantially reduce the energy use and nutrient loss associated with conventional techniques. As a result, aseptically packaged products retain more nutritional value and exhibit more natural texture, color, and taste.\nThe rapid heating and cooling steps of the aseptic process typically take place while the product is being piped from one vessel to another. A portion of the pipe is heated and, as the product passes through the heated portion of the pipe, heat is transferred to the product. The product within the pipe is under pressure, typically 50 to 100 psi, which serves to move the product through the pipe and also prevent vaporization of the water component of the product within superheated sections. Subsequent operations, such as packaging, normally require lower pressures, typically up to 15 psi. Thus, the product must go though a pressure reduction step before packaging.\nMany designs of back-pressure or pressure reduction valves exist for single-phase fluids. However, these valves cause significant shear damage to sensitive fluids and particulate-containing liquids. Moreover, the valves can become clogged with particulates, which can damage the valve and potentially cause a rupture.\nAccordingly, there is a need to provide a pressure reduction valve that can be used to transfer a continuous flow of particulate-containing material from a higher pressure environment to a lower pressure condition without significant damage to the particulates contained in the material."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor assembly processes, and more particularly, to a method and apparatus for picking up a semiconductor chip by blowing air through air holes in a dicing tape to separate the semiconductor chip from the dicing tape and then picking up the semiconductor chip. A method of removing a semiconductor chip from a dicing tape and a method for forming a perforated dicing tape are also provided.\n2. Description of the Related Art\nIn semiconductor assembly processes, dividing a wafer into semiconductor chips may be termed “sawing” or “dicing.” After the wafer has been diced or sawed into individual semiconductor chips, each of the semiconductor chips may be separated from the wafer and attached to a lead frame. To perform a dicing process, dicing tape may be attached to the back of the wafer. A needle pick-up method is conventionally used to separate a semiconductor chip from the dicing tape. However, because the needle applies physical stress to the semiconductor chip, the needle may damage the chip as it forcibly pushes the semiconductor chip off the dicing tape.\nReferring to FIG. 1, in a conventional needle pick-up method, a wafer 20 may be mounted to a dicing tape 10 (P1). Next, the wafer 20 may be divided into separate chips 30 by sawing or dicing the wafer (P2). UV radiation 40 may be applied to weaken the adhesion of the dicing tape 10 (P3). The chips 30 may then be removed from the dicing tape 10 (e.g., “picked up”) using a needle 50 that may have one or more pins 55. For example, the needle 50 may be positioned under the dicing tape 10 (and the semiconductor chip 30), thereby pushing the semiconductor chip 30 off the dicing tape 10. The pushed-up semiconductor chip 30 may then be picked up by a suction collection member 60 positioned over the chip 30. The semiconductor chip 30 may then be attached to a mounting pad of a lead frame (not shown).\nWhen the needle pick-up method described above is used to remove the semiconductor chip 30 from the dicing tape 10, the sharp end of the needle pins 55 used to penetrate the dicing tape 10 and force the semiconductor chip 30 off of the dicing tape 10 may apply physical stress to the semiconductor chip 30. This stress on the semiconductor chip 30 may physically damage the semiconductor chip 30, such as by cracking the semiconductor chip 30 or breaking the edges of the semiconductor chip 30. In addition, because semiconductor chip thickness has been reduced so that thin semiconductor packages that have a thickness of several tens of μm or multi chip packages that have a thickness of several hundred μm or less may be formed, the conventional needle pick-up method may cause severe damage to such a thin chip. Because a damaged semiconductor chip 30 is not easily separated from the dicing tape 10 and leads to defective units, manufacturing costs may be increased."}
{"text": "Plastic optical materials are lightweight, hardly fragile and tintable compared to inorganic optical materials. Various plastic materials, such as resins, are currently used in optical materials and are gradually required to have better physical properties.\nPolythiourethane optical resins produced using polythiol compounds and isocyanate compounds are widely used as optical lens materials due to their excellent optical properties, including high transparency, Abbe number, transmittance and tensile strength. However, optical materials produced by curing resin compositions including polythiol compounds and general-purpose isocyanate compounds suffer from frequent nonuniformity polymerization or whitening, which worsens the optical properties of the optical resins. General-purpose isocyanate compounds and polythiol compounds as main components of resin compositions for thiourethane lenses are prone to nonuniformity polymerization, whitening and clouding depending on their miscibility. In contrast, optical resins produced by heat curing of highly miscible isocyanate compounds and polythiol compounds do not substantially suffer from the problems of nonuniformity polymerization and whitening even when the polythiol compounds are not specially treated. Such isocyanate compounds include, for example, 3,8-bis(isocyanatomethyl)tricyclo[5.2.1.02.6]decane, 3,9-bis(isocyanatomethyl)tricyclo[5.2.1.02.6]decane, 4,8-bis(isocyanatomethyl)tricyclo[5.2.1.02.6]decane, 4,9-bis(isocyanatomethyl)tricyclo[5.2.1.02.6]decane, 2,5-bis(isocyanatomethyl)bicyclo[2.2.1]heptane, and 2,6-bis(isocyanatomethyl)bicyclo[2.2.1]heptanes. However, the isocyanate compounds are difficult to prepare, which incurs considerable preparation costs. As a consequence, the use of the isocyanate compounds inevitably increases the production cost of thiourethane lenses. Meanwhile, optical resins obtained by curing inexpensive general-purpose isocyanates and polythiol compounds suffer from frequent nonuniformity polymerization or whitening. Such isocyanate compounds include, for example, isophorone diisocyanate, dicyclohexylmethane-4,4-diisocyanate (H12MDI), and 1,6-hexamethylenediisocyanate. Particularly, tape whitening and bubbling occur in some cases. These phenomena adversely affect the performance of optical materials and are causes of high defective proportion and low lens quality."}
{"text": "This invention relates in general to data communication, and more particularly to a loop integrity test device and method for digital subscriber line (xDSL) communication.\nDigital subscriber line (xDSL) communication can provide a high bandwidth data path supported by the twisted pair wiring infrastructure of the public switched telephone network (PSTN). xDSL technologies support data service simultaneously with traditional telephone service using a separation technique. Suitable xDSL technologies include asymmetric digital subscriber line (ADSL), rate adaptable digital subscriber line (RADSL), symmetric digital subscriber line (SDSL), high-speed digital subscriber line (HDSL) and very high-speed digital subscriber line (VDSL).\nIn general, an xDSL communication link is established across the local loop between customer premises equipment (CPE) and the local loop termination point (e.g., central office or remote terminal). The local loop termination point typically has a xDSL access multiplexer (DSLAM) that handles the xDSL link on the network side, and the CPE typically includes one or more xDSL termination units (XTU) that handle the xDSL link on the customer premises side. One architecture for xDSL communication is disclosed in U.S. Pat. No. 5,668,857 Sep. 16, 1997, and entitled xe2x80x9cCommunication Server Apparatus and Method.xe2x80x9d\nIn order to support xDSL communication, the twisted pair line on the local loop between the customer premises and the loop termination point must meet certain physical characteristics. If it does not, an xDSL communication link can not be successfully established. Consequently, it can be desirable to allow the xDSL capability of the local loop to be quickly verified at the same time that xDSL service is installed at the customer premises. Further, it can be desirable to monitor the xDSL capability of the local loop on an ongoing and continuous basis.\nIn accordance with the present invention, a loop integrity test device and method for digital subscriber line communication are disclosed that provide advantages over conventional xDSL communication devices and systems.\nAccording to one aspect of the present invention, a digital subscriber line (xDSL) communication system allows xDSL communication across a local loop. The system includes a local loop termination point and customer premises equipment connected to a twisted pair telephone line. Loop integrity test devices for xDSL communication are located at the local loop termination point and the customer premises equipment and are coupled to the telephone line. The loop integrity test devices are respectively operable to transmit test signatures across the telephone line, to receive and evaluate test signatures from the telephone line, and to indicate whether the telephone line can support xDSL communication based upon evaluation of test signatures.\nAccording to another aspect of the present invention, a loop integrity test device is disclosed that includes a line interface unit operable to connect to and interface with a twisted pair telephone line. The test device also includes an output device operable to indicate a pass state and a fail state. Further, the test device includes a control unit having a generation unit and an evaluation unit. The generation unit is operable to transmit a test signature across the telephone line. The evaluation unit is operable to receive and evaluate a test signature to determine whether the telephone line can support xDSL communication. The control unit is then operable to direct the output device to indicate a pass state if the telephone line can support DSL communication and to indicate a fail state if the telephone line can not support xDSL communication.\nA technical advantage of the present invention is that the integrity of the local loop can be verified at the time of installation of high speed xDSL data service at a customer premises. In particular, the present loop integrity test device provides an immediate indication of whether the loop can support xDSL communication when it is installed.\nAnother technical advantage of the present invention is an improvement in the reliability and monitoring of the xDSL communication system. The test devices implement a continuous integrity check of the local loop that allows the telephone company or other service provider to be alerted to a problem on the local loop. This notification may allow the service provider to correct the problem before a degradation of service is noticed by the customer.\nOther technical advantages should be readily apparent to one skilled in the art from the following figures, description, and claims."}
{"text": "1. Field of the Invention\nThe present invention relates to a film image information reading apparatus (ordinarily, referred to as a film scanner) for reading image information of a film by a photoelectric conversion element and, more particularly, to the adjustment of a gain of an output level of the photoelectric conversion element when the image of a negative film is read.\n2. Related Background Art\nHitherto, there has generally been known a color scanner (color image information reading apparatus) in which various color information of a color document is separated into three primary colors of, e.g., R (red), G (green), and B (blue), and three kinds of image signals corresponding to the respective color separation images are obtained from the photoelectric conversion element, and the color information of the color document is reproduced by synthesizing those three kinds of image signals.\nHowever, if a film such as photographic film, microfilm, or the like is used as a document, in general, the density range of the film is fairly wide to be above 3. In particular, if a self scan type sensor such as a CCD (charge coupled device) or the like is used as the photoelectric conversion element of the color scanner, the density range of the film is fairly wider than the dynamic range of the sensor. Therefore, the image signal read by such a sensor is largely adversely influenced and a problem occurs. Practically speaking, there is such a problem that the gradation of the portion having a high or low density of the film image is not reproduced.\nOn the other hand, if a color negative film is used as a document, as shown by film characteristic curves of FIG. 2, the density of each layer of R, G, and B in the unexposed portions is considerably large. Therefore, it will be appreciated that the color negative film is constituted in such a manner that the negative image is reproduced on the orange base portion. In other words, in the case of the negative film the mixed image information consisting of the true image information and the color information of the film base portion which is unconcerned with the image information is read as the color image signal. Therefore, the dynamic range of the sensor is not effectively used. In addition, in the case of the negative film, the negative image is formed. When the negative image is reversed to the positive image, the density of the base portion is reduced, so that the density information of the base portion is quite the meaningless image information. Further, the orange components of the base portion of the document slightly differs in dependence on the film maker, kind of film, or developing state of the film, causing a problem when the accurate color is reproduced."}
{"text": "This invention relates to a charged particle beam apparatus comprising means for generating a beam of charged particles, and a lens system including a magnetic lens for focussing the beam.\nThe beam of charged particles may be used to effect a process at a surface of a target. For example, in the field of semiconductor technology it is known to use a charged particle beam apparatus in the form of a so-called electron beam column to direct an electron beam towards a semiconductor wafer coated with an electron sensitive resist and to write a predetermined pattern in the resist by computer control of the beam. After exposure the resist is developed to form a masking layer which can be used subsequently in the processing of the semiconductor wafer.\nA conventional electron beam column produces a Gaussian, round beam spot which is scanned across the target to write the desired pattern. Although this exposure technique has the advantage of Providing optimum pattern flexibility it suffers from the drawback of being relatively slow.\nShorter exposure times can be achieved by using an electron beam column capable of producing a larger, square-shaped beam spot of fixed size. In this case the pattern is written by moving the shaped beam in discrete steps across the target.\nRecent developments in the technology of electron beam columns have made it possible to shorten exposure times still further by using a rectangular-shaped beam spot whose size and shape can be altered to fit the various areas of the pattern to be exposed. This technique is known as variable-shaped beam imaging. The shaping variations are performed while the beam is stepped from one position to the next.\nThe lens system of a typical variable-shaped electron beam column basically comprises five magnetic lenses. In this context reference is invited to U.S. Pat. No. 4,243,866 and more particularly to the description with reference to FIG. 6. While the electron-beam column described there comprises various other elements, the basic lens constituents are as follows. Firstly, after the electron source there are two magnetic lenses associated with the beam shaping facility. Then, moving along the length of the column away from the electron source, there is a pair of demagnifying lenses and finally a so-called projection lens. Essentially these lenses are all arranged in series in the sense that, travelling the length of the column, the electron beam experiences the field of each lens in turn.\nIn the present context a magnetic lens normally comprises an electrically conductive coil which is enclosed within a soft ferromagnetic member forming the pole pieces of the magnet. The coil is disposed in the vicinity of the electron beam path so that the beam of electrons can be focussed by the action of the magnetic field produced by the passage of direct current through the coil. The magnetic field strength produced by the magnet is dependent on the magnitude of the current flowing in the coil so that the focussing effect (i.e. the strength) of the lens can be altered simply by increasing or decreasing the current. It is noted here that magnetic focussing is inevitably accompanied by rotation of the beam and the extent of this rotation is dependent on the strength of the magnetic lens.\nHowever, there are occasions when the focussing of the electron beam column has to be altered within a limited range while the machine is actually operating. This may be necessary, for example, when there is a change in the height of the target on which the electron beam is impinging or when the size of the beam spot (and hence the beam current) is to be altered. In order not to erode the advantage of the fast exposure times which are capable with the variable-shaped electron beam column the refocussing should be completed in a time period similar to that taken to change the spot size, typically tens of nanoseconds. Unfortunately the time taken to change the strength of a magnetic lens by altering the current flow is considerably greater than this because of the inductance of the coil. Therefore, although the variable-shaped electron beam column is capable of fast exposure times, the number of wafers which can be processed per unit time by the machine can be limited by the relatively long times necessary for adequate refocussing."}
{"text": "Cholinergic receptors normally bind the endogenous neurotransmitter acetylcholine (ACh), thereby triggering the opening of ion channels. ACh receptors in the mammalian central nervous system can be divided into muscarinic (mAChR) and nicotinic (nAChR) subtypes based on the agonist activities of muscarine and nicotine, respectively. The nicotinic acetylcholine receptors are ligand-gated ion-channels containing five subunits. Members of the nAChR subunit gene family have been divided into two groups based on their amino acid sequences; one group containing so-called (3 subunits, and a second group containing a subunits. Three kinds of a subunits, α7, α8 and α9, have been shown to form functional receptors when expressed alone and thus are presumed to form homooligomeric pentameric receptors.\nAn allosteric transition state model of the nAChR has been developed that involves at least a resting state, an activated state and a “desensitized” closed channel state, a process by which receptors become insensitive to the agonist. Different nAChR ligands can stabilize the conformational state of a receptor to which they preferentially bind. For example, the agonists ACh and (−)-nicotine respectively stabilize the active and desensitized states.\nChanges of the activity of nicotinic receptors has been implicated in a number of diseases. Some of these, for example myasthenia gravis and ADNFLE (autosomal dominant nocturnal front lobe epilepsy) are associated with reductions in the activity of nicotinic transmission either because of a decrease in receptor number or increased desensitization.\nReductions in nicotinic receptors have also been hypothesized to mediate cognitive deficits seen in diseases such as Alzheimer's disease and schizophrenia.\nThe effects of nicotine from tobacco are also mediated by nicotinic receptors. and since the effect of nicotine is to stabilize receptors in a desensitized state, an increased activity of nicotinic receptors may reduce the desire to smoke.\nCompounds which bind nAChRs have been suggested for the treatment of a range of clinical psychiatric and neurological disorders involving cognitive deficits in attention, alertness, and memory. These may include those that may benefit from selective enhancement in cholinergic transmission such as attention deficit, chronic psychotic disorders, jetlag, selected pain syndromes, and smoking cessation; and those thought to involve reduced cholinergic function such as, neurodegenerative disorders, central inflammatory or autoimmune diseases, brain trauma and cerebral vascular disease. Modulation of α7 nicotinic receptor activity is expected to be beneficial in a number of diseases including Alzheimer's disease, Mild Cognitive Impairment (MCI) and related syndromes, Lewy Body Dementia, vascular dementia, Multiple Sclerosis, post-encephalitic dementias, Parkinson's disease, Huntington's disease, Attention Deficit Hyperactivity Disorder (ADHD), schizophrenia, bipolar mood disorder, schizoaffective disorder, Tourette's syndrome, and brain trauma.\nHowever, treatment with nicotinic receptor agonists which act at the same site as ACh is problematic because ACh not only activates, but also blocks receptor activity through processes which include desensitization and uncompetitive blockade. Furthermore, prolonged activation appears to induce a long-lasting inactivation. Therefore, agonists of ACh can be expected to reduce activity as well as enhance it.\nAt nicotinic receptors in general, and of particular note at the α7-nicotinic receptor, desensitization limits the duration of action of an applied agonist."}
{"text": "1. Field of the Invention\nThe present invention relates to a lift device and, more particularly, to a lift device including a plurality of inflatable bladders.\n2. Description of Related Art\nIndividuals who lack a capacity to stand or move themselves from the prone or supine position into a seated or standing position must rely on outside assistance to achieve the desired position. Currently, there are a variety of methods to assist such an individual including assistance by attendants, assistance by loved ones or significant others, assistance by Trained Nurses, or assistance by means of mechanical devices. Each of these methods, however, suffers from various drawbacks.\nFor instance, if the assistant is one person, such as an attendant or nurse, then that person must be strong enough to navigate both himself and a patient. The individual must be experienced and know the proper procedure for moving a patient from a down position. There are frequent injuries to both patient and assistant when an assistant attempts to move a fallen patient. According to the Bureau of Labor Statistics, nursing aides, orderlies and attendants reported more musculoskeletal disorders requiring time off from work than any other occupation. Many of these disorders were due to overexertion related to lifting or moving patients.\nCurrent mechanical lift devices also suffer from a variety of drawbacks. These devices are very bulky, heavy and difficult to maneuver and store. Furthermore, current mechanical lifts can create bruising and stress marks due to the focused forces caused by the straps when a patient's weight is not evenly distributed. An example of such a mechanical lift device can be found in U.S. Pat. No. 6,857,144 to Huang.\nInflatable lift devices also exist in the prior art, however, these devices also suffer from various drawbacks. U.S. Pat. No. 5,669,086 to Garman, for instance, discloses a mobile lifting apparatus that includes a stack of inflatable bags for lifting a patient off the ground. However, this device requires a ramp for positioning the patient on the device. Therefore, this device cannot be positioned beneath a patient without moving the patient. Furthermore, this device merely lifts a patient off the ground and does not include inflatable side rails, armrests, and a backrest that helps the patient from a position on the ground to a seated position and, finally, to a standing position.\nAccordingly, a need exists for a device that easily lifts a patient who has fallen without the need for an assistant to exert strenuous force. A further need exists for an inflatable lift device that provides support for a patient's arms and back while slowly lifting the patient to a seated or standing position."}
{"text": "1. Field of the Invention\nThis invention relates to a method and apparatus for displaying fluorescence information, wherein excitation light is irradiated to a region of interest in a living body, intrinsic fluorescence produced by an intrinsic dye in the living body is detected, and information in accordance with characteristics of the intrinsic fluorescence is displayed.\n2. Description of the Related Art\nThere have heretofore been proposed fluorescence information displaying techniques utilizing characteristics such that, in cases where excitation light having wavelengths falling within an excitation wavelength range for an intrinsic dye in a living body is irradiated to the living body, an intensity of fluorescence produced by the intrinsic dye in the living body varies for normal tissues and diseased tissues. With the proposed fluorescence information displaying techniques, excitation light having predetermined wavelengths is irradiated to a region of interest in a living body, the fluorescence produced by an intrinsic dye in the living body is detected, and the location and the infiltration range of diseased tissues are displayed as a fluorescence image.\nOrdinarily, when excitation light is irradiated to a region of interest in a living body, the fluorescence having a high intensity is produced by normal tissues, and the fluorescence having a low intensity is produced by diseased tissues. Therefore, by measurement of the fluorescence intensity, a judgment as to the state of a disease is capable of being made.\nBasically, apparatuses for displaying fluorescence information comprise excitation light irradiating means for irradiating excitation light, which has wavelengths falling within an excitation wavelength range for an intrinsic dye in a living body, to the living body, imaging means for detecting fluorescence produced by the intrinsic dye in the living body and forming a fluorescence image of the living body, and image displaying means for receiving the output from the imaging means and displaying the fluorescence image. In many cases, the apparatuses for displaying fluorescence information take on the form built in endoscopes, which are inserted into the body cavities, colposcopes, operating microscopes, or the like.\nHowever, the aforesaid apparatuses for displaying fluorescence information have the problems described below. Specifically, since a region in a living body has protrusions and recesses, the distance between the light source of the excitation light irradiating means and the region of interest in the living body is not uniform. Therefore, ordinarily, the irradiance of the excitation light at the living body portion, which is exposed to the excitation light, is non-uniform. The intensity of fluorescence is approximately in proportion to the irradiance of the excitation light, and the irradiance of the excitation light at the portion, which is exposed to the excitation light, is in inverse proportion to the square of the distance between the light source of the excitation light irradiating means and the portion, which is exposed to the excitation light. Accordingly, the problems occur in that diseased tissues, which are located close to the light source, produce the fluorescence having a higher intensity than the intensity of the fluorescence produced by normal tissues, which are located remote from the light source. The problems also occur in that the intensity of the fluorescence from normal tissues, which are located at a position inclined with respect to the excitation light, becomes markedly low. Thus if the irradiance of the excitation light is non-uniform, the intensity of the fluorescence will vary in accordance with the level of the irradiance of the excitation light, and therefore an error will often be made in the judgment of the state of a disease.\nA fluorescence imaging technique has been proposed in, for example, xe2x80x9cFluorescence Imaging of Early Lung Cancer,xe2x80x9d Annual International Conference of the IEEE Engineering in Medicine and Biology Society, Vol. 12, No. 3, 1990. With the proposed technique, the fluorescence, which is produced by an intrinsic dye in an region of interest in a living body when the region of interest is exposed to excitation light, is separated into an intensity of the green wavelength region (hereinbelow referred to as the xe2x80x9cgreen region intensity Gxe2x80x9d) and an intensity of the red wavelength region (hereinbelow referred to as the xe2x80x9cred region intensity Rxe2x80x9d). An image operation is then performed in accordance with division of the red region intensity R and the green region intensity G by each other, and the results of the division are displayed. The proposed technique utilizes the findings in that the spectrum of the fluorescence produced by normal tissues is different from the spectrum of the fluorescence produced by diseased tissues. Specifically, when the spectrum of the fluorescence, which is produced by the intrinsic dye at normal tissues in the living body, and the spectrum of the fluorescence, which is produced by the intrinsic dye at diseased tissues in the living body, are compared with each other, in particular, the intensity of the green region of the spectrum obtained from the diseased tissues is markedly lower than the intensity of the green region of the spectrum obtained from the normal tissues. Therefore, the degree of reduction in the intensity of the green region intensity G of the fluorescence, which is produced from the diseased tissues, as compared with the intensity of the green region intensity G of the fluorescence produced from the normal tissues, is markedly higher than the degree of reduction in the intensity of the red region intensity R of the fluorescence, which is produced from the diseased tissues, as compared with the intensity of the red region intensity R of the fluorescence produced from the normal tissues. Therefore, only the fluorescence from the diseased tissues can be specifically extracted by the division of R/G and can be displayed as an image.\nSpecifically, with the proposed technique, the term of the fluorescence intensity depending upon the distance between the excitation light source and the region of interest in the living body and the distance between the region of interest in the living body and the fluorescence receiving means is canceled, and information reflecting only the difference in fluorescence spectrum pattern is obtained.\nHowever, heretofore, research has not been conducted sufficiently with respect to a combination of detection wavelengths, at which the difference between the pattern of a fluorescence spectrum obtained from normal tissues and the pattern of a fluorescence spectrum obtained from diseased tissues occurs markedly. Therefore, there have been the problems in that a desirable combination of detection wavelengths cannot be presented with numerical values.\nThe primary object of the present invention is to provide a method of displaying fluorescence information, wherein fluorescence components having two different wavelengths are extracted from fluorescence, the two different wavelengths having been specified with numerical values as an appropriate combination of detection wavelengths, at which a difference between a pattern of a fluorescence spectrum obtained from normal tissues and a pattern of a fluorescence spectrum obtained from diseased tissues occurs markedly, light intensities of the extracted fluorescence components are detected, and information in accordance with a ratio between the light intensities of the extracted fluorescence components is displayed with a high reliability.\nAnother object of the present invention is to provide a method of displaying fluorescence information, wherein fluorescence components having wavelengths falling within a certain wavelength region are extracted from fluorescence, the certain wavelength region having been specified with numerical values as an appropriate detection wavelength region, at which a difference between a fluorescence intensity of fluorescence produced from normal tissues and a fluorescence intensity of fluorescence produced from diseased tissues occurs markedly, a light intensity of the extracted fluorescence components is detected, and information in accordance with the detected light intensity is displayed with a high reliability.\nThe specific object of the present invention is to provide an apparatus for carrying out the method of displaying fluorescence information.\nThe present invention provides a first method of displaying fluorescence information, comprising the steps of:\ni) irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) detecting light intensity of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a predetermined wavelength region containing 480 nm,\niii) detecting light intensity of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a predetermined wavelength region containing either one of 630 nm and 700 nm, and\niv) displaying information in accordance with a ratio between the two detected light intensities.\nThe present invention also provides a first apparatus for displaying fluorescence information, comprising:\ni) excitation light irradiating means for irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) first fluorescence intensity detecting means for detecting light intensity of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a predetermined wavelength region containing 480 nm,\niii) second fluorescence intensity detecting means for detecting light intensity of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a predetermined wavelength region containing either one of 630 nm and 700 nm, and\niv) displaying means for displaying information in accordance with a ratio between the light intensity, which has been detected by the first fluorescence intensity detecting means, and the light intensity, which has been detected by the second fluorescence intensity detecting means.\nThe first apparatus for displaying fluorescence information in accordance with the present invention should preferably be modified such that the first fluorescence intensity detecting means comprises:\nfirst wavelength selecting means for selecting the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region containing 480 nm, and\nfirst light intensity detecting means for detecting the light intensity of the fluorescence components having been selected by the first wavelength selecting means, and\nthe second fluorescence intensity detecting means comprises:\nsecond wavelength selecting means for selecting the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region containing either one of 630 nm and 700 nm, and\nsecond light intensity detecting means for detecting the light intensity of the fluorescence components having been selected by the second wavelength selecting means.\nAlso, the first apparatus for displaying fluorescence information in accordance with the present invention should more preferably be modified such that the first wavelength selecting means selects the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region of 480 nmxc2x1at most 70 nm, and\nthe second wavelength selecting means selects the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region of either one of 630 nmxc2x1at most 70 nm and 700 nmxc2x1at most 70 nm.\nFurther, in the first apparatus for displaying fluorescence information in accordance with the present invention, the excitation light should preferably have wavelengths falling within the range of 380 nm to 420 nm, which range is apart from the characteristic peak light intensity of the fluorescence produced from normal tissues of the living body. Furthermore, in the first apparatus for displaying fluorescence information in accordance with the present invention, the excitation light irradiating means should preferably be a GaN type of semiconductor laser.\nThe first apparatus for displaying fluorescence information in accordance with the present invention may be constituted so as to two-dimensionally detect a fluorescence image. Alternatively, the first apparatus for displaying fluorescence information in accordance with the present invention may be constituted so as to detect the fluorescence intensity with respect to each point at a site in the living body.\nAlso, in the first apparatus for displaying fluorescence information in accordance with the present invention, by way of example, each of the first wavelength selecting means and the second wavelength selecting means may be constituted so as to extract the fluorescence components, which have wavelengths falling within the predetermined wavelength region, with a dichroic mirror, or the like. Alternatively, each of the first wavelength selecting means and the second wavelength selecting means may be constituted so as to extract the fluorescence components, which have wavelengths falling within the predetermined wavelength region, in a time division mode by changing over an optical filter, or the like. As another alternative, in cases where a fluorescence image is to be detected two-dimensionally, each of the first wavelength selecting means and the second wavelength selecting means may be constituted so as to extract the fluorescence components, which have wavelengths falling within the predetermined wavelength region, by utilizing-a mosaic filter formed by connecting optical filters with one another in a mosaic pattern.\nFurther, in the first apparatus for displaying fluorescence information in accordance with the present invention, the displaying means may employ one of various displaying techniques. For example, the displaying means may be constituted so as to display the ratio between the light intensity of the fluorescence components, which have wavelengths falling within the predetermined wavelength region containing 480 nm, and the light intensity of the fluorescence components, which have wavelengths falling within the predetermined wavelength region containing either one of 630 nm and 700 nm, on a monitor, with a printer, or the like. Alternatively, the displaying means may be constituted so as to alter the tint or the luminance of the displayed color in accordance with the ratio between the two light intensities.\nThe present invention further provides a second method of displaying fluorescence information, comprising the steps of:\ni) irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) detecting light intensity B of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm, and\niii) displaying information in accordance with the light intensity B.\nThe present invention still further provides a third method of displaying fluorescence information, comprising the steps of:\ni) irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) detecting light intensity W of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within an entire measurement wavelength region,\niii) detecting at least one light intensity selected from among:\nlight intensity Bxe2x80x2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm,\nlight intensity R1 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 630 nmxc2x1at most 70 nm and at least containing 600 nm to 630 nm, and\nlight intensity R2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm, and\niv) displaying information in accordance with a ratio between the at least one selected light intensity and the light intensity W.\nThe present invention also provides a fourth method of displaying fluorescence information, comprising the steps of:\ni) irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) detecting light intensity B of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm,\niii) detecting light intensity W of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within an entire measurement wavelength region,\niv) detecting at least one light intensity selected from among:\nlight intensity Bxe2x80x2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm,\nlight intensity R1 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 630 nmxc2x1at most 70 nm and at least containing 600 nm to 630 nm, and\nlight intensity R2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm, and\nv) displaying information in accordance with the light intensity B and a ratio between the at least one light intensity, which is selected from among the light intensities Bxe2x80x2, R1, and R2, and the light intensity W.\nIn the third and fourth methods of displaying fluorescence information in accordance with the present invention, in cases where the light intensity Bxe2x80x2 is detected as the at least one light intensity, the light intensity Bxe2x80x2 may be the light intensity of the fluorescence components having wavelengths falling within the wavelength region identical with the wavelength region of the fluorescence components, whose light intensity is detected as the light intensity B. Alternatively, the light intensity Bxe2x80x2 may be the light intensity of the fluorescence components having wavelengths falling within the wavelength region different from the wavelength region of the fluorescence components, whose light intensity is detected as the light intensity B. In the former cases, instead of the light intensity Bxe2x80x2 being detected besides the light intensity B, the light intensity B may be employed as the light intensity Bxe2x80x2. The term xe2x80x9cdetecting light intensity Bxe2x80x2 xe2x80x9d as used herein also includes the cases where the light intensity B is employed as the light intensity Bxe2x80x2.\nThe present invention further provides a second apparatus for displaying fluorescence information, comprising:\ni) excitation light irradiating means for irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) first light intensity detecting means for detecting light intensity B of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm, and\niii) fluorescence information displaying means for displaying information in accordance with the light intensity B having been detected by the first light intensity detecting means.\nThe present invention still further provides a third apparatus for displaying fluorescence information, comprising:\ni) excitation light irradiating means for irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) second light intensity detecting means for detecting light intensity W of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within an entire measurement wavelength region,\niii) at least one light intensity detecting means selected from among:\nthird light intensity detecting means for detecting light intensity Bxe2x80x2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm,\nfourth light intensity detecting means for detecting light intensity R1 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 630 nmxc2x1at most 70 nm and at least containing 600 nm to 630 nm, and\nfifth light intensity detecting means for detecting light intensity R2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm, and\niv) fluorescence information displaying means for displaying information in accordance with a ratio between the light intensity, which has been detected by the at least one selected light intensity detecting means, and the light intensity W.\nThe present invention also provides a fourth apparatus for displaying fluorescence information, comprising:\ni) excitation light irradiating means for irradiating excitation light to a measuring site in a living body, the excitation light causing the measuring site to produce fluorescence,\nii) first light intensity detecting means for detecting light intensity B of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm,\niii) second light intensity detecting means for detecting light intensity W of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within an entire measurement wavelength region,\niv) at least one light intensity detecting means selected from among:\nthird light intensity detecting means for detecting light intensity Bxe2x80x2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm,\nfourth light intensity detecting means for detecting light intensity R1 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 630 nmxc2x1at most 70 nm and at least containing 600 nm to 630 nm, and\nfifth light intensity detecting means for detecting light intensity R2 of fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm, and\nv) fluorescence information displaying means for displaying information in accordance with the light intensity B and a ratio between the light intensity, which has been detected by the at least one selected light intensity detecting means, and the light intensity W.\nIn the third and fourth apparatuses for displaying fluorescence information in accordance with the present invention, in cases where the third light intensity detecting means is provided as the at least one light intensity detecting means, the third light intensity detecting means may be identical with the first light intensity detecting means. Alternatively, the third light intensity detecting means may be different from the first light intensity detecting means. In the former cases, instead of the third light intensity detecting means being provided besides the first light intensity detecting means, the first light intensity detecting means may be utilized as the third light intensity detecting means. The term xe2x80x9ccomprising third light intensity detecting meansxe2x80x9d as used herein also includes the cases where the first light intensity detecting means is utilized as the third light intensity detecting means.\nThe second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention should preferably be modified such that the light intensity detecting means comprises an image sensor for two-dimensionally detecting the fluorescence produced by the measuring site and forming a fluorescence image, and wavelength selecting means for selecting the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a desired wavelength region.\nAlso, the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention may be modified such that the light intensity detecting means is provided with fluorescence acquiring means for acquiring the fluorescence, which is produced by a single point at a site in the living body.\nFurther, in the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, the excitation light should preferably have wavelengths falling within the range of 380 nm to 420 nm, which range is apart from the characteristic peak light intensity of the fluorescence produced from normal tissues of the living body. Furthermore, in the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, the excitation light irradiating means should preferably be a GaN type of semiconductor laser.\nFurthermore, in the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, the fluorescence information displaying means may employ one of various displaying techniques. For example, in cases where the light intensity of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within a predetermined wavelength region containing 480 nm, and the light intensity of the fluorescence components of the fluorescence, which fluorescence components have wavelengths falling within the entire measurement wavelength region, are detected and the information in accordance with the ratio between the two detected light intensities are displayed, the fluorescence information displaying means may be constituted so as to display the ratio between the two detected light intensities on a monitor, with a printer, or the like. Alternatively, the fluorescence information displaying means may be constituted so as to alter the tint or the luminance of the displayed color in accordance with the ratio between the two detected light intensities.\nThe first method of displaying fluorescence information and the first apparatus for displaying fluorescence information in accordance with the present invention have the effects described below.\nWhen the excitation light is irradiated to living body tissues, the living body tissues are excited by the excitation light and produce the fluorescence having a spectrum illustrated in FIG. 5. It is assumed that the thus produced fluorescence results from superposition of the fluorescence produced by various kinds of intrinsic dyes in the living body, such as FAD, collagen, fibronectin, and porphyrin. FIG. 5 shows typical fluorescence spectra of the fluorescence produced from normal tissues and the fluorescence produced from diseased tissues, which fluorescence spectra have been measured by the inventors.\nAs illustrated in FIG. 5, the level and the pattern of the spectrum of the fluorescence vary for the normal tissues and the diseased tissues. The level of the fluorescence produced from the normal tissues is high as a whole, and the level of the fluorescence produced from the diseased tissues is low as a whole. The fluorescence spectrum obtained from the normal tissues has a peak spectral intensity at a region in the vicinity of 480 nm, which region is a blue region. The fluorescence spectrum obtained from the diseased tissues has a peak spectral intensity at a region in the vicinity of 630 nm, which region is a red region, and a peak spectral intensity at a region in the vicinity of 700 nm.\nFIG. 6 is a graph showing a distribution of ratios of spectral intensities at respective wavelengths of a fluorescence spectrum, which is obtained from normal tissues, to a spectral intensity of an entire measurement wavelength region width, which spectral intensity is taken as 1, and a distribution of ratios of spectral intensities at respective wavelengths of a fluorescence spectrum, which is obtained from diseased tissues, to a spectral intensity of the entire measurement wavelength region width, which spectral intensity is taken as 1. As illustrated in FIG. 6, the distributions of the spectral intensity ratios more clearly manifest the difference between the pattern of the fluorescence spectrum of the fluorescence produced from the normal tissues and the pattern of the fluorescence spectrum of the fluorescence produced from the diseased tissues.\nIn accordance with FIG. 5 and FIG. 6, the inventors studied about an appropriate combination of detection wavelengths, at which the difference between the pattern of the fluorescence spectrum obtained from the normal tissues and the pattern of the fluorescence spectrum obtained from the diseased tissues occurs markedly.\nAs a result, it has been found that, in cases where the light intensity at the wavelength region in the vicinity of 480 nm and the light intensity at the wavelength region in the vicinity of 630 nm or in the vicinity of 700 nm are detected, the ratio between the two detected light intensities markedly represents the difference between the pattern of the fluorescence spectrum obtained from the normal tissues and the pattern of the fluorescence spectrum obtained from the diseased tissues.\nSpecifically, the light intensity at the predetermined wavelength region containing 480 nm, at which a high light intensity is obtained characteristically from the normal tissues, and the light intensity at the predetermined wavelength region containing 630 nm or 700 nm, at which a high light intensity is obtained characteristically from the diseased tissues, are detected from the fluorescence having been detected from the measuring site, whose tissue state is unknown. Also, the information in accordance with the ratio between the two detected light intensities is displayed. In such cases, a person, who sees the displayed information, is capable of presuming whether the tissues at the measuring site are the normal tissues or the diseased tissues.\nAs described above, with the first method of displaying fluorescence information and the first apparatus for displaying fluorescence information in accordance with the present invention, the excitation light is irradiated to the measuring site in the living body, the excitation light causing the measuring site to produce the fluorescence. Also, the light intensity of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region containing 480 nm, is detected. Further, the light intensity of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region containing either one of 630 nm and 700 nm. Thereafter, the information in accordance with the ratio between the two detected light intensities is displayed. Therefore, with the first method of displaying fluorescence information and the first apparatus for displaying fluorescence information in accordance with the present invention, the information having enhanced reliability is capable of being displayed.\nIn the first apparatus for displaying fluorescence information in accordance with the present invention, the first fluorescence intensity detecting means may comprise the first wavelength selecting means for selecting the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region containing 480 nm, and the first light intensity detecting means for detecting the light intensity of the fluorescence components having been selected by the first wavelength selecting means. Also, the second fluorescence intensity detecting means may comprise the second wavelength selecting means for selecting the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region containing either one of 630 nm and 700 nm, and the second light intensity detecting means for detecting the light intensity of the fluorescence components having been selected by the second wavelength selecting means. With the thus modified first apparatus for displaying fluorescence information in accordance with the present invention, the light intensity at each of the predetermined wavelength regions is capable of being detected easily.\nAlso, from the results of analysis made on the fluorescence spectra illustrated in FIG. 5 and FIG. 6, it has been found that, if the width of the predetermined wavelength region containing 480 nm is broader than 480 nmxc2x170 nm, the ratio of the light intensity at the predetermined wavelength region detected from the diseased tissues to the light intensity at the predetermined wavelength region detected from the normal tissues will become high. Further, it has been found that, if the width of the predetermined wavelength region containing 630 nm is broader than 630 nm xc2x170 nm, or if the width of the predetermined wavelength region containing 700 nm is broader than 700 nmxc2x170 nm, the ratio of the light intensity at the predetermined wavelength region detected from the diseased tissues to the light intensity at the predetermined wavelength region detected from the normal tissues will become low. Therefore, with the first apparatus for displaying fluorescence information in accordance with the present invention, wherein the first wavelength selecting means selects the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region of 480 nm xc2x1at most 70 nm, and the second wavelength selecting means selects the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the predetermined wavelength region of either one of 630 nmxc2x1at most 70 nm and 700 nmxc2x1at most 70 nm, the light intensity ratio having reliability enhanced even further is capable of being obtained.\nFurther, with the first apparatus for displaying fluorescence information in accordance with the present invention, wherein the excitation light has wavelengths falling within the range of 380 nm to 420 nm that is apart from the region in the vicinity of 480 nm, at which the light intensity of the fluorescence produced from the normal tissues becomes characteristically high, the fluorescence having the fluorescence spectrum of a desirable pattern is capable of being produced, and therefore the reliability of the displayed information is capable of being enhanced. Furthermore, with the first apparatus for displaying fluorescence information in accordance with the present invention, wherein the GaN type of semiconductor laser is employed as the excitation light irradiating means, the size of the apparatus is capable of being kept small, and the cost of the apparatus is capable of being kept low.\nThe second, third, and fourth methods of displaying fluorescence information and the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention have the effects described below.\nAs illustrated in FIG. 5, the level and the pattern of the spectrum of the fluorescence vary for the normal tissues and the diseased tissues. The level of the intrinsic fluorescence produced from the normal tissues is high as a whole, and the level of the intrinsic fluorescence produced from the diseased tissues is low as a whole. The fluorescence spectrum obtained from the normal tissues has a peak spectral intensity at a region in the vicinity of 480 nm, which region is the blue region. However, at the region in the vicinity of 480 nm, only slight fluorescence is produced from the diseased tissues. From FIG. 5, it can be found that the wavelengths, at which the difference between the spectral intensity of the fluorescence spectrum obtained from the normal tissues and the spectral intensity of the fluorescence spectrum obtained from the diseased tissues occurs markedly, are the wavelengths in the vicinity of 480 nm, and the wavelength of 480 nm is desirable as the center wavelength of the detection wavelengths.\nAlso, in experiments made by the inventors and with respect to a plurality of patients, fluorescence spectra of the fluorescence produced from diseased tissues and the fluorescence produced from normal tissues located in the vicinity of the diseased tissues were detected. Thereafter, from each of the fluorescence spectra, light intensities at various measurement wavelength regions having different widths and having their centers at 480 nm were calculated. With respect to each of the widths of the measurement wavelength regions, the ratio of the light intensity of the fluorescence produced from the diseased tissues to the light intensity of the fluorescence produced from the normal tissues was calculated. From the calculations, the results shown in FIG. 11 were obtained.\nFrom the results shown in FIG. 11, it can be found that, in cases where the width of the measurement wavelength region is at most 100 nm, i.e. in cases where the measurement wavelength region is 480 nm at most 50 nm, little change occurs in the ratio of the light intensity of the fluorescence produced from the diseased tissues to the light intensity of the fluorescence produced from the normal tissues. Also, it can be found that, in cases where the width of the measurement wavelength region is 150 nm, the ratio of the light intensity of the fluorescence produced from the diseased tissues to the light intensity of the fluorescence produced from the normal tissues becomes high and close to 1, and the contrast between the fluorescence produced from the normal tissues and the fluorescence produced from the diseased tissues becomes low. Therefore, it can be found desirable to extract the fluorescence components having wavelengths falling within the wavelength region width of at most 140 nm around the wavelength of 480 nm, i.e. the wavelengths falling within the wavelength region of 480 nmxc2x1at most 70 nm, at which wavelength region little increase occurs with the ratio of the light intensity of the fluorescence produced from the diseased tissues to the light intensity of the fluorescence produced from the normal tissues.\nFurther, from the results of experiments, it has been found that the wavelength corresponding to the peak spectral intensity in the fluorescence spectrum often varies slightly for different persons and different measuring sites in a single person, and therefore the light intensity at the wavelength region, which at least contains 450 nm to 480 nm, should be detected.\nFIG. 6 shows the distribution of ratios of spectral intensities at respective wavelengths of the fluorescence spectrum, which is obtained from the normal tissues, to the spectral intensity of the entire measurement wavelength region width, which spectral intensity is taken as 1, and the distribution of ratios of spectral intensities at respective wavelengths of the fluorescence spectrum, which is obtained from the diseased tissues, to the spectral intensity of the entire measurement wavelength region width, which spectral intensity is taken as 1. In cases where the ratios of the spectral intensities at respective wavelengths of the fluorescence spectrum to the spectral intensity of the entire measurement wavelength region width are calculated, adverse effects of a difference in measurement conditions, such as a measurement distance, can be eliminated. Therefore, the distributions of the spectral intensity ratios manifest the difference between the pattern of the fluorescence spectrum of the fluorescence produced from the normal tissues and the pattern of the fluorescence spectrum of the fluorescence produced from the diseased tissues. The fluorescence spectrum obtained from the normal tissues has the peak spectral intensity at the region in the vicinity of 480 nm. The fluorescence spectrum obtained from the diseased tissues has the peak spectral intensity at the region in the vicinity of 630 nm, and the peak spectral intensity at the region in the vicinity of 700 nm.\nIn accordance with FIG. 6, the inventors studied about measurement wavelengths and measurement wavelength region width, at which the difference between the pattern of the fluorescence spectrum obtained from the normal tissues and the pattern of the fluorescence spectrum obtained from the diseased tissues occurs markedly. The study was made in the same manner as that described above and by setting the center wavelengths at 480 nm, 630 nm, and 700 nm. FIG. 12 shows the results obtained from experiments, in which the center wavelength was set at 480 nm, the measurement wavelength region width was set at various different values, the ratio of the light intensity of each measurement wavelength region width to the light intensity of the entire measurement wavelength region width was calculated with respect to each of the fluorescence produced from the normal tissues and the fluorescence produced from the diseased tissues, and the ratio of the thus calculated light intensity ratio for the diseased tissues to the thus calculated light intensity ratio for the normal tissues was calculated.\nAs a result, from the study made with the spectral intensity ratio distribution, as in the study made with the spectral intensity distribution, in cases where the center wavelength is 480 nm, it has been found desirable to extract the fluorescence components having wavelengths falling within the wavelength region width of at most 140 nm around the wavelength of 480 nm, i.e. the wavelengths falling within the wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm, at which wavelength region little increase occurs with the ratio of the light intensity of the fluorescence produced from the diseased tissues to the light intensity of the fluorescence produced from the normal tissues and the contrast between the fluorescence produced from the normal tissues and the fluorescence produced from the diseased tissues does not become low.\nAlso, though not shown, in cases where the center wavelength is set at 630 nm, it has been found desirable to extract the fluorescence components having wavelengths falling within the wavelength region of 630 nmxc2x1at most 70 nm and at least containing 600 nm to 630 nm. Further, in cases where the center wavelength is set at 700 nm, it has been found desirable to extract the fluorescence components having wavelengths falling within the wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm.\nSpecifically, the fluorescence components having the wavelengths falling within the wavelength region described above are extracted from the fluorescence having been detected from the measuring site, whose tissue state is unknown, and the light intensity of the extracted fluorescence components is detected. Also, the information in accordance with the detected light intensity is displayed. In such cases, a person, who sees the displayed information, is capable of accurately presuming whether the tissues at the measuring site are the normal tissues or the diseased tissues.\nAs described above, with the second method of displaying fluorescence information and the second apparatus for displaying fluorescence information in accordance with the present invention, the excitation light is irradiated to the measuring site in the living body, the excitation light causing the measuring site to produce the fluorescence. Also, the light intensity B of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 480 nm xc2x1at most 70 nm and at least containing 450 nm to 480 nm, is detected. Further, the information in accordance with the light intensity B is displayed. Therefore, with the second method of displaying fluorescence information and the second apparatus for displaying fluorescence information in accordance with the present invention, the information having enhanced reliability is capable of being displayed.\nWith the third method of displaying fluorescence information and the third apparatus for displaying fluorescence information in accordance with the present invention, the excitation light is irradiated to the measuring site in the living body, the excitation light causing the measuring site to produce the fluorescence. Also, the light intensity W of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the entire measurement wavelength region, is detected. Further, at least one light intensity is detected, the at least one light intensity being selected from among the light intensity Bxe2x80x2 of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm, the light intensity R1 of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 630 nm xc2x1at most 70 nm and at least containing 600 nm to 630 nm, and the light intensity R2 of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm. Furthermore, the information in accordance with the ratio between the at least one selected light intensity and the light intensity W is displayed. Therefore, with the third method of displaying fluorescence information and the third apparatus for displaying fluorescence information in accordance with the present invention, adverse effects of fluctuations in spectral intensity due to fluctuations in measurement conditions, such as the measurement distance and the measurement angle, are capable of being reduced, and the information having enhanced reliability, which represents the feature of the spectrum pattern, is capable of being displayed.\nWith the fourth method of displaying fluorescence information and the fourth apparatus for displaying fluorescence information in accordance with the present invention, the excitation light is irradiated to the measuring site in the living body, the excitation light causing the measuring site to produce the fluorescence. Also, the light intensity B of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm, is detected. Further, the light intensity W of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the entire measurement wavelength region, is detected. Furthermore, at least one light intensity is detected, the at least one light intensity being selected from among the light intensity Bxe2x80x2 of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 480 nmxc2x1at most 70 nm and at least containing 450 nm to 480 nm, the light intensity R1 of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 630 nm xc2x1at most 70 nm and at least containing 600 nm to 630 nm, and the light intensity R2 of the fluorescence components of the fluorescence produced by the measuring site, which fluorescence components have wavelengths falling within the wavelength region of 700 nmxc2x1at most 70 nm and at least containing 700 nm to 710 nm. Also, the information in accordance with the light intensity B and the ratio between the at least one light intensity, which is selected from among the light intensities Bxe2x80x2, R1, and R2, and the light intensity W is displayed. Therefore, with the fourth method of displaying fluorescence information and the fourth apparatus for displaying fluorescence information in accordance with the present invention, the information having enhanced reliability, which represents the characteristics of the spectral intensity and the spectrum pattern, is capable of being displayed.\nWith the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, wherein each of the light intensity detecting means is provided with the image sensor for two-dimensionally detecting the fluorescence produced by the measuring site and forming a fluorescence image, the information concerning the fluorescence over a wide range is capable of being displayed quickly.\nWith the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, wherein each of the light intensity detecting means is provided with the fluorescence acquiring means for acquiring the fluorescence, which is produced by a single point at a site in the living body, the information concerning the fluorescence produced from a desired point at the site is capable of being displayed.\nWith the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, wherein the excitation light has the wavelengths falling within the range of 380 nm to 420 nm that is apart from the region in the vicinity of 480 nm, at which region the light intensity of the fluorescence produced from normal tissues takes the characteristically large value, the fluorescence with the fluorescence spectrum in a desirable pattern is capable of being produced, and the reliability of the displayed information is capable of being enhanced. Furthermore, with the second, third, and fourth apparatuses for displaying fluorescence information in accordance with the present invention, wherein the GaN type of semiconductor laser is employed as the excitation light irradiating means, the size of the apparatus is capable of being kept small, and the cost of the apparatus is capable of being kept low."}
{"text": "This invention relates to a tape cassette having a case which is smaller, especially in thickness, than a conventional tape cassette and which can still accommodate a magnetic tape of the same width as can a conventional tape cassette, and to an adaptor for mounting such a tape cassette to a conventional tape recorder.\nAmong cassette tapes, the most widespread now in use are called compact cassettes, which contain a magnetic tape having a width of 3.8 mm and have a standardized shape and size, i.e. a length of 100 mm, a width of 64 mm and a thickness of 8.5 mm (FIG. 8).\nA handy tape recorder for this type of compact cassette is rather small and lightweight. In fact, the latest model is only slightly larger than the cassette itself. In other words, further reduction in size of the hardware part of the tape recorder is almost impossible.\nOn the other hand, a micro-cassette is known which is extremely small (50 mm in length, 33 mm in width and 8 mm in thickness). But since its sound information recording capacity is smaller than a compact cassette and it cannot feed the tape stably due to its fragile build, its sound quality is considerably poor. Thus, tape recorders for micro-cassettes are not popular in spite of their compact size.\nFor the above reasons, it is essential, in attempting to realize a reduction in size of a tape cassette, not to impair the sound information recording capacity, not to impair the sound quality and to keep a sufficient strength of the cassette itself.\nAlso, it is desirable that such a small and thin tape cassette be capable of being recorded with and being replayed on a conventional compact cassette tape recorder.\nJapanese Unexamined Patent Publication No. 56-98757 discloses a compact cassette in which a tape cassette is inserted in a case to protect its magnetic tape during transportation and which can be mounted on a tape recorder with the tape cassette received in the case.\nAlso, Japanese Examined Patent Publication No. 57-58752 discloses a solid type tape cassette having a solid interior except at some portions. It has an increased strength."}
{"text": "Field of the Invention\nThe present invention relates to a decorative cover for containers, and more particularly but not by way of limitation, to a decorative cover for flower pots and to methods for producing such decorative cover wherein the decorative cover contains design indicia to personalize or enhance the visual aesthetic effect of the decorative cover."}
{"text": "1. Field of the Invention\nThe present invention relates to a spark plug used for ignition in an internal combustion engine.\n2. Description of the Related Art\nFIG. 9 shows a conventional spark plug 300 used for ignition in an internal combustion engine, such as an automobile gasoline engine. The spark plug 300 is mounted on a cylinder head SH of an engine by a male-threaded portion 301a formed on an outer circumferential surface of a metallic shell 301. When the spark plug 300 is mounted on the cylinder head SH, a spark discharge gap g defined by a ground electrode 304 and a center electrode 303 is positioned within a combustion chamber BR and is adapted to ignite a fuel-air mixture. A hexagonal portion 305 (tool engagement portion) is formed on an outer circumferential surface of the metallic shell 301 and is adapted to tighten the male-threaded portion 301a by a tightening tool. The metallic shell 301 assumes a substantially cylindrical shape having a bore 306 for receiving an insulator 302 and is conventionally manufactured by cold plastic processing and machining. In many spark plugs, in order to improve manufacturing efficiency, a schematic profile and the bore 306 are formed by die forging, and a final profile including the male-threaded portion 301a is finished by machining. Since the metallic shell 301 has a thin-walled portion, the dimensions of the metallic shell 301 must be designed in consideration of a material flow during die forging; otherwise, a working defect is likely to occur.\nWith a recent tendency toward complication of engine head structure, space allocated around a valve for installation of the spark plug 300 is decreasing. Thus, the hexagonal portion 305 needs to be reduced to increase space for use on the head side. However, reducing the size of the hexagonal portion 305 causes the following problems.\n(1) To prevent an excessive reduction in the wall thickness of the hexagonal portion 305 in association with a reduction in the size of the hexagonal portion 305, a diameter D1 of a portion (hereinafter referred to as xe2x80x9ca major-bore portionxe2x80x9d) 306a of the shell bore 306 must be reduced. Also, the outside diameter of the insulator 302 must be reduced accordingly. However, when a diameter D2 of a portion (hereinafter referred to as xe2x80x9can intermediate-bore portionxe2x80x9d) of the shell bore 306 corresponding to the male-threaded portion 301a is reduced, a forging punch becomes excessively thin, when the intermediate-bore portion 306b is formed by forging and thus may be damaged or may cause a working defect when a large working load is applied thereon. This problem arises particularly in the case when the male-threaded portion 301a has a long screw reach.\n(2) A portion of the insulator 302 positioned within the major-bore portion 306a is formed into a flange portion 302e. When the metallic shell 306 is swaged onto the insulator 302, the flange portion 302e bears a swaging force. A metallic terminal 313 and a center electrode 303 are connected by a glass seal portion 315. In the step of the glass seal portion 315, the flange portion 302e bears a pressing force. In particular, the center electrode 303, a material powder of the glass seal portion 315, and the metallic terminal 313 are disposed within a through-hole formed in the insulator 302. Then the insulator 302 is inserted into a bore formed in a seat die such that the flange portion 302e rests on an inner seat portion formed on the wall of the bore. In this state, the entire insulator 302 is heated to a temperature equal to or higher than a glass softening point, and the metallic terminal 313 is pressed inwardly in the axial direction to press the material powder with the center electrode 303, thereby forming the glass seal portion 315. During this pressing process, the flange portion 302e bears a pressing force.\nIf the outside diameter of the insulator 302 is too small to meet the demand described above in (1), manufacturing the insulator 302 becomes very difficult. Therefore, there is a certain limit to a reduction in the outside diameter of the insulator 302. As the size of the hexagonal portion 305 is reduced, the diameter D1 of the major-bore portion 306a is reduced accordingly. Thus, the diameter of the flange portion 302e, which is accommodated within the major-bore portion 306a, is also reduced. Because of a reduction in the size of the hexagonal portion 305, the diameter of the flange portion 302e must be reduced because there is a certain limit to a reduction in the diameter of a portion of the insulator 302 other than the flange portion 302e (for example, a portion of the insulator 302 positioned within the intermediate-bore portion 306b; hereinafter referred to as xe2x80x9can intermediate-trunk portion 302axe2x80x9d). As a result, the amount of a projection of the flange portion 302e decreases, causing, for example, a decrease in the area of contact between the flange portion 302e and the seat portion of the seat die used in the step of forming the glass seal portion 315. Consequently, a load concentration causes breakage of the seat die or galling of the insulator 302 and the seat die.\n(3) If the diameter of the intermediate-trunk portion 302a of the insulator 302 is reduced to meet the demand described above in (2), and also the diameter D2 of the intermediate-bore portion 306b of the metallic shell 306 is set to a rather large value to attain favorable workability during the process in (1), a gap is likely to be formed between the intermediate-bore portion 306b and the intermediate-trunk portion 302a of the insulator 302. The presence of this gap tends to cause an eccentric disposition of the insulator 302 within the metallic shell 301, potentially causing an impairment in spark plug performance (for example, lateral sparking).\nAn object of the present invention is to provide a spark plug capable of increasing the degree of freedom with respect to space around a cylinder head on which the spark plug is mounted, through reduction in the size of a tool engagement portion, such as a hexagonal portion, and capable of implementing the following:\n(1) in spite of a reduction in the size of the tool engagement portion, a metallic shell can be manufactured efficiently and at high yield;\n(2) during formation of a conductive glass seal layer or a resistor by use of a seat die, breakage or galling of the seat die is less likely to occur; and\n(3) during incorporation of an insulator into the metallic shell, an eccentric disposition of the insulator within the metallic shell is less likely to occur.\nTo achieve the above object, the present invention provides a spark plug including an elongate center electrode, an insulator enclosing the center electrode, a metallic shell having open opposite ends and enclosing the insulator, the metallic shell having a male-threaded portion formed on a front-side outer circumferential surface of the metallic shell, and a tool engagement portion formed on the outer circumferential surface of the metallic shell at a rear side with respect to the male-threaded portion, the tool engagement portion projecting circumferentially outwardly, and a ground electrode connected to the metallic shell and defining a spark discharge gap in cooperation with the center electrode.\nIn the specification, the term xe2x80x9cfrontxe2x80x9d refers to a spark discharge gap side with respect to an axial direction of the center electrode, and the term xe2x80x9crearxe2x80x9d refers to a side opposite the front side.\nThe insulator has a stepped annular insulator-side engagement portion for engaging with an annular shell-side engagement portion projected inwardly from a portion of an inner surface of the metallic shell corresponding to the male-threaded portion, and |Axe2x88x92E|xe2x89xa61.5 mm, and 0.4xe2x89xa6(D2/E)2xe2x89xa60.6, where A is a dimension of the tool engagement portion represented by a diameter of an inscribed circle of a cross-sectional outline of the tool engagement portion, E is an effective diameter of the male-threaded portion, and D2 is an inner diameter of an intermediate-bore portion of the metallic shell located on a rear side with respect to the shell-side engagement portion.\nAccording to the above-described structure, the dimension A of the tool engagement portion (for example, a hexagonal portion) is reduced with respect to the effective diameter E of the male-threaded portion such that |Axe2x88x92E| becomes not greater than 1.5 mm. Thus, the degree of freedom with respect to space around a cylinder head on which the spark plug is mounted can be increased. Even when available space around a valve for installation of the spark plug decreases due to complication of cylinder head structure, the spark plug can be easily mounted on the cylinder head. Although the outside diameter of the insulator decreases in association with a reduction in the size of the tool engagement portion, so long as 0.4xe2x89xa6(D2/E)2xe2x89xa60.6, the wall thickness of the male-threaded portion of the metallic shell falls within an appropriate range. Thus, during forging of the metallic shell, a forging punch is less susceptible to breakage and is less likely to cause a working defect. That is, the problem described previously in (1) is solved, and the metallic shell can be manufactured efficiently and at high yield.\nMore particularly, (D2/E)2 represents the ratio of the cross-sectional area of the intermediate-bore portion having the diameter D2 xe2x80x9cxcfx80(D2/2)2xe2x80x9d to the cross-sectional area of the male-threaded portion having the effective diameter E xe2x80x9cxcfx80(E/2)2.xe2x80x9d The smaller the value (D2/E)2 (i.e., the more the effective diameter E of the male-threaded portion increases with respect to the diameter D2 of the intermediate-bore portion), the greater the wall thickness of the male-threaded portion. When (D2/E)2 is less than 0.4, the wall thickness of the male-threaded portion becomes excessively large, causing an insufficient diameter of the intermediate-bore portion. As a result, when the intermediate-bore portion is to be formed through cold working, such as forging, a forging punch to be used becomes excessively thin and thus may be damaged or may cause a working defect when a large working load is imposed thereon. When (D2/E)2 is in excess of 0.6, the wall thickness of the male-threaded portion becomes excessively thin. As a result, formation of the male-threaded portion through cold working becomes difficult, and the formed male-threaded portion suffers insufficient strength. More preferably, (D2/E)2 ranges from 0.45 to 0.55.\nA flange portion may be formed on the outer circumferential surface of the insulator on the rear side with respect to the stepped portion. In this case, preferably, d2/d1 is not greater than 0.75, where d1 is the outside diameter of the flange portion, and d2 is the outside diameter of an intermediate-trunk portion extending between the flange portion and the stepped portion. As mentioned previously in (2), in the case of reducing the outside dimension A of the tool engagement portion such that |Axe2x88x92E| is not greater than 1.5 mm, if the outside diameter of the intermediate-trunk portion becomes excessively small, manufacture of the insulator becomes very difficult. Also, a reduction in the size of the tool engagement portion unavoidably requires a reduction in the outside diameter of the flange portion. In other words, the diameter ratio d2/d1 between the intermediate-trunk portion and the flange portion tends to increase. As d2/d1 increases, the amount of projection of the flange portion from the outer circumferential surface of the intermediate-trunk portion decreases. As a result, as mentioned previously, the step of forming a glass seal portion is likely to involve breakage of a seat die or galling between the insulator and the seat die. Through employment of a d2/d1 of not greater than 0.7, the amount of projection of the flange portion becomes sufficiently large, thereby effectively preventing the above-mentioned problem associated with a reduction in the size of the tool engagement portion; i.e., solving the problem described previously in (2). More preferably, d2/d1 is not greater than 0.65. However, d2/d1 is excessively small, the intermediate-trunk portion becomes too thin for manufacture of the insulator. Therefore, in order to avoid such a problem, the value d2/d1 must be adjusted as adequate.\nAs mentioned previously in (3), if the diameter of the intermediate-bore portion is set to a rather large value in order to attain favorable workability of the metallic shell while the diameter of the intermediate-trunk portion of the insulator is decreased in association with a reduction in the size of the tool engagement portion, a gap is likely to be formed between the intermediate-bore portion of the metallic shell and the intermediate-trunk portion of the insulator. In the case of formation of such a gap, preferably, an eccentricity preventive portion is provided substantially concentrically with the intermediate-bore portion and the intermediate-trunk portion in such a manner as to partially fill the gap. In the step of incorporating the insulator into the metallic shell, the eccentricity preventive portion restricts lateral movement of the insulator; i.e., an eccentric disposition of the insulator within the metallic shell, thereby solving the problem described previously in (3).\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory only and are not restrictive of the invention, as claimed."}
{"text": "The use of soap dispensers for hand washing has now become widely known and accepted. Typically, such soap dispensers dispense a quantity of soap which is then worked into a lather by the user when combined with water on the hands. Recently, there has been a general acceptance of foam soap delivery systems. In such systems, liquid soap is combined with air, typically under force or pressure, and then driven through a mesh, screen or porous passage to finish or homogenize the soap into a uniform stable composition. In some systems, a mixing chamber is employed prior to the porous structure or passage in order to prepare a prefoam of randomly sized and spaced bubbles. The mixing chamber and porous passage are generally presented at a foamer head that is immediately adjacent the drive mechanism for the liquid and air. In general, these drive systems are typically pistons within cylinders or pumps to achieve the pressurization and drive of the liquid soap and air.\nIn the past, little attention has been given to the development of foamer heads that are adaptable for use in systems where the liquid soap source is remote from the dispensing head. Indeed, the prior art foamer heads have typically been of a rudimentary nature, with little regard for the specifics of the design or the configuration of the constituent elements. While the prior art foamer heads have generally been of a satisfactory nature, little attention has been given to the efficacy of soap foam generation to achieve a desired uniformity and integrity of the resulting foam. Moreover, where foam soap is to be dispensed from an area remote from the liquid soap source, the prior art has generally taught the generation of the foam close to the liquid soap source, with its subsequent delivery to a dispensing head remote from that source. However, such systems have generally proven to be problematic. It has been found that foam is difficult to drive for any distance through a conduit. Breakdown of the foam occurs, resulting in reduced volumes of soap being dispensed on each dispensing cycle, and with the ultimate dispensing of liquid soap globules. It has also been found that such remote delivery systems have resulted in extremely low output volumes on subsequent dispensing operations, and even total failures to dispense when the period of time between dispensing operations has been sufficient to allow the soap foam within the conduit to fully breakdown. Other problems have been evidenced with a “wet” foam output on subsequent dispensing operations, resulting from the breakdown of foam in the conduit into a liquid form.\nIn systems where the dispensing head is remote from the point of foam generation, it has been found that the liquid and/or air cylinders of this system have required careful design to ensure sufficient “suck-back” force on the return stroke of the dispensing operation to draw residual foam back away from the dispensing head to preclude drips and the like.\nThe remote dispensing heads referenced herein are typically present in what are referred to as counter-mount systems, in which the soap reservoir is maintained beneath the counter and the dispensing head is above the counter, the two being interconnected by conduits that are three or more feet in length. The problems of foam breakdown and suck-back failure are characteristic of such systems.\nThere is a need in the art for an improved soap foam generator, adaptable for use in any of a variety of delivery systems, and particularly in remote dispensing systems, such as counter-mount systems, in which the air and liquid sources are remote from the dispensing head."}
{"text": "The present invention relates to vehicle internal safety belts, and more particularly to a vehicle internal safety belt for folding one's head to the head-rest of the vehicle internal seat, which can also be alternatively used for holding one's legs together.\nConventionally, a vehicle internal seat is generally equipped with a safety belt for fastening on one's waist to firmly secure a person to a seat. However, even if a person is fastened to a vehicle internal seat by a safety belt, a person's head is still not protected and may be shaken to hit a window glass or some part inside a vehicle during vehicle collision or sudden braking of the vehicle. Further, there is no means available in a vehicle seat to bind one's legs together while one is sleeping in a vehicle internal seat. Because it is ungraceful to unintentionally open the two legs apart while a women is sleeping in a vehicle internal seat, a woman may be constantly in a state of anxiety while sleeping in a vehicle internal seat."}
{"text": "The present invention relates to azlactone activated polyalkylene oxides (PAO's) having improved hydrolytic stability, and to water-soluble polyalkylene oxide conjugates prepared therefrom.\nThe conjugation of water-soluble polyalkylene oxides with useful molecules such as proteins and polypeptides is well known. The coupling of peptides and polypeptides to polyethylene glycol (PEG) and similar water-soluble polyalkylene oxides is disclosed by U.S. Pat. No. 4,179,337 to Davis et al.\nDavis et al. discloses that physiologically active polypeptides modified with PEG exhibit dramatically reduced immunogenicity and antigenicity. Also, the polyalkylene oxide conjugates, when injected into a living organism, have been shown to remain in the bloodstream considerably longer than the corresponding native proteins. Examples of such therapeutic protein conjugates include tissue plasminogen activator, insulin, interleukin II and hemoglobin. In addition, PAO's have also been conjugated to oligonucleotides. See, for example, U.S. Pat. No. 4,904,582.\nTo conjugate polyalkylene oxides, the hydroxyl end-groups of the polymer must first be converted into reactive functional groups. This process is frequently referred to as \"activation\" and the product is called an \"activated polyalkylene oxide.\"\nUntil recently, covalent attachment of the polyalkylene oxide to an appropriate nucleophile was effected by activated polyalkylene oxides such as polyalkylene oxide succinoyl-N-hydroxysuccinimide ester, as disclosed by Abuchowski et al., Cancer Biochem. Biophys., 7, 175-86 (1984). This polyalkylene oxide derivative is desirable because it is reactive under mild conditions.\nA shortcoming associated with this derivative, however, is the fact that it is relatively hydrolytically unstable when no nucleophile is present. Recently, in U.S. Pat. No. 5,122,614, polyalkylene oxide-N-succinimide carbonates were disclosed having improved hydrolytic stability over the polyalkylene oxide succinoyl succinates. Even so, the pH conditions necessary to deprotonate the .epsilon.-NH.sub.2 groups of polypeptide lysines for conjugation subject the activated polyalkylene oxide to hydrolysis. This does not affect the reaction end product, other than to reduce its yield. While reduced yields ordinarily affect product cost, the hydrolysis becomes even more costly for several reasons. Firstly, reaction mixtures cannot be prepared significantly in advance. Additional purification of the end product is required to remove the hydrolytic degradation products. Furthermore, the reduction in yield is compensated for by increasing the amount of activated polyalkylene oxide starting material. This increases the viscosity of the reaction mixture, thereby further increasing the processing cost, and potentially interferes with downstream purification of the polymer and conjugate.\nA need exists, therefore, for polyalkylene oxides that are unreactive towards weak nucleophiles such as water but react readily with stronger nucleophiles such as polypeptides. While azlactones have been reported to react readily with amines and less readily with water, azlactone activated PAO's are unreported. Unsaturated azlactones, in particular, are not altered by long contact with water. See, Carter, Organic Reactions, Vol. III (Adams, ed., John Wiley & Sons, New York 1946) pp. 198-239. The disclosed unsaturated azlactones have the following structures ##STR1## in which R.sub.1 and R.sub.2 are independently selected from hydrogen, phenyl rings and lower alkyl moieties. R is the residue of an .alpha.-acyl amino acid."}
{"text": "Though the current invention is also applicable to lithographic processes as used in printing industry, the current invention is mainly concerned with lithography as applied in the production of electronic and electric devices, sensors, and the like. The characteristic length scale, i.e. the design rule of the relevant lithographic pattern is usually in the micron to submicron region. The manufacturing process of these devices is currently almost entirely dominated by methods based on optical or e-beam lithography.\nSince the emergence of integrated circuits (ICs) and micromechanical devices, optical lithography has been crucial for the purpose of their mass production: Its convenience, parallel operation and resolution has created a huge market. Fabrication of devices with ever smaller dimensions, necessary to satisfy the demands of storage and computation, becomes increasingly problematic with visible light as processes steadily reach fundamental limits, predominantly set by diffraction. This realization triggered intense research in UV, X-ray, e-beam, and Scanning Probe (SP)\nlithography. These methods deliver high resolution with varying success and their economics remain, at best, uncertain. Reasons for these uncertainties include limitations due to wavelength dependent phenomena, the slow writing speeds of e-beam and SP lithographies, and challenges in finding appropriate resists and masks.\nA separate and related limitation of current lithographies is the complexity of processes required for modifying or altering a layer at regions defined by the desired lithographic pattern to be written; lithography today relies on bulk transfer of reactive material from the liquid or gas phase, or from a plasma using masks to protect the other regions of the substrate.\nAn alternative approach to lithography has been published by A. Kumar and G. M. Whitesides in : Appl. Phys. Lett. 1993, 63, 2002-2004. In this process, known as microcontact stamp lithography, stamps are fabricated by casting a replica in poly(dimethylsiloxane) (PDMS) of a master with a negative of the desired pattern. In the one known example of microcontact stamping used for lithography, the PDMS stamp is inked with an alkanethiol, hexadecanethiol, and transferred to gold by transient contact between the stamp and the gold substrate. The thiol covalently binds to the gold in an autophobic reaction: the thiol modifies the wettability of the gold substrate, preventing spread of the bulk liquid phase and thus confining the transferred monolayer to the region of contact between the raised regions of the elastomeric stamp and the substrate. The presence of these thiols allows subsequent lithographic processing of the gold using a cyanide/oxygen etch that selectively removes gold not protected by a monolayer of alkanethiol. Although this system allows reproductions of features in gold down to one micron, its scope remains limited to a special subset of useful materials, i.e. thiols and gold.\nIn all these lithographic processes, a number of agents are used which pose a considerable potential threat to the environment. Huge efforts are therefore taken to sufficiently protect workers within IC factories and the population living in their surrounding from a possible exposure to these agents. Immense resources are also invested for safely depositing the waste and remains of the IC manufacturing.\nThe object of the present invention is to avoid limitations in the type of reactions, materials transferred and range of substrates useful in contact based lithography and to demonstrate new strategies of direct processing of substrates without the intervening application of resists.\nIt is a further object of the inventions to improve the known lithographic processes by providing a method for modifying a surface at precisely determined positions. While being more convenient to apply, the new process should at least have a resolution comparable to state-of-the-art lithography. Another object of the invention is to restrict the amount of agents and material used in the lithographic process, in particular, hazardous waste generated by this process should be minimized.\nAs mentioned above, lithographic processes are also found in the field of printing. In particular, color printing is found to face similar problems concerning alignment during several separated printing steps as encountered in lithography for semiconductor devices. In four color printing the colors cyan, magenta, yellow and black are applied to four separate printing rollers that transfer the lithographic pattern to the paper in four separate printing steps. The alignment among those four printing steps limits the maximal resolution and quality of color images especially on cheap and low quality paper. Printing can either be done with poor alignment (.about.0.1 mm) and bad image quality on cheap paper or with better alignment (.about.0.01 mm) and better image quality on high quality paper. The rastering dimension has to be adapted to the alignment and consequently is &gt;0.1 mm (&lt;300 dpi) for low quality paper and&gt;0.01 mm (&lt;3000 dpi) for the paper with the highest quality. Photographic reproduction, on the other hand, can be carried out with grain sizes of the order of microns (&lt;0.001 mm). Thus, high quality printing requires expensive equipment able to maintain tolerances of 10 microns over several meters.\nIt is therefore seen as another object of the invention to improve the conventional color printing procedure, in particular to facilitate the alignment procedure necessary to print several colors."}
{"text": "Exocrine glands are glands that release a secretion external to or at the surface of an organ by means of a canal or duct. Examples of exocrine glands include, among others, the mammary glands, prostate, liver, gall bladder, pancreas, kidneys, sweat glands, and salivary glands. Cancers of exocrine glands pose a major health problem, frequently resulting in death. Currently, cancers of the breast and prostate are among the leading causes of death among women and men, respectively.\nThe mature human breast comprises from six to nine major ducts, which emanate from the nipple, serially branch into ducts and terminate in lobuloalveolar structures (Russo et al., Lab. Invest. 62(3): 244–278 (1990)). This branching network of ducts is composed of epithelial cells in a supporting matrix of connective tissue and endothelial cells.\nTissues removed from the human female breast during surgery and autopsy have been examined in numerous studies directed to the nature and site of origin of neoplastic growth. Subgross sampling and histological confirmation have enabled pathological characterization of entire breasts, leading to the postulation of the existence of four major possible sites of origin of mammary carcinomas, namely ducts, terminal ducts, ductules, and acini (Russo et al., supra). Ductal origin is supported by the presence of more extensive epithelial proliferations, which are presumed to be preneoplastic, in surgically removed cancerous breasts as compared to nonmalignant breasts removed during autopsies (Russo et al., supra).\nWith a cumulative lifetime risk of a woman developing breast cancer estimated to be 1 in 9, there is an urgent need to develop therapeutic methods of treatment that are more effective, less invasive and accompanied by fewer side effects and prophylactic methods of treatment that are more effective than increased and intensified physical monitoring and far less extreme than radical mastectomy. In spite of the recent discovery of the heritable breast cancer susceptibility loci, BRCA1 (Miki et al., Science 266: 66–71 (1994)) and BRCA2, and other cancer susceptibility loci, and the increasing ability of physicians to identify women with elevated breast cancer risk, prophylactic methods are still currently limited to physical monitoring and prophylactic mastectomy.\nIn view of the above, it is an object of the present invention to provide a method of locally treating an exocrine gland prophylactically for disease. It is another object of the present invention to provide a method of locally treating an exocrine gland, in particular the mammary gland, prophylactically for cancer. Another object of the present invention is to provide a method of locally treating an exocrine gland therapeutically for disease. Yet another object of the present invention is to provide a method of locally treating a mammary gland therapeutically for cancer. Still yet another object of the present invention is to provide a method of locally treating a mammary gland both therapeutically and prophylactically for cancer. These and other objects and advantages, as well as additional inventive features, will become apparent from the detailed description provided herein."}
{"text": "The present invention relates to Caterpillar sidebooms used for pipelaying, and more particularly pertains to methods and apparatus for converting old, discontinued, mechanically-operated, low-drive, Caterpillar sidebooms into hydraulically-operated sidebooms, with minimal retrofitting.\nIt is well known in the art that crawler or tractor-type vehicles having an integrated, rigid maneuverable boom disposed on a side thereof are commonly used for pipelaying operations, i.e., for raising, carrying, and lowering heavy pipe. Referred to as \"sidebooms,\" such vehicles must be capable of safely handling heavy pipe; indeed, there are sideboom models that have been constructed for handling pipe up to 200,000 pounds. Such sidebooms were mechanically-operated in their first incarnation and are now readily available from original equipment manufacturers in hydraulically-operated models.\nAn early development in the mechanical sideboom art is disclosed by Butterfield et al. in U.S. Pat. No. 3,785,503, wherein a plurality of planetary gears and concomitant shafts are used to drive the winch drums for each of the boom and the load. Primarily due to its inherent complicated, high-maintenance gearing system, no embodiments of the Butterfield sideboom was ever commercialized. A significant improvement in the pipelaying sideboom art would be a sideboom that eliminates all gears to accomplish transmission of power, i.e., that eliminates all mechanical connections between the engine and the winch system.\nOther early developments in the sideboom art, albeit not applicable to pipelaying operations, per se, are taught in U.S. Pat. Nos. 2,909,290; and 3,329,283. More particularly, Nichols, in U.S. Pat. No. 2,909,290, teaches a farm tractor-mounted sideboom intended for lifting light loads typical on the farm. Of course, pipelaying operations demand sidebooms with lifting capacities up to 200,000 pounds. Similarly, Wade, in U.S. Pat. No. 3,329,283, teaches a snap mount sideboom configured to be foldable for reducing its prerequisite overhead clearance. Half of the Wade boom's height may be reduced and the folded boom portion secured to the tractor's side by using integral hooks and cable. Especially in view of the Wade sideboom having no counterweight, it should be evident to those skilled in the art that both Wade and Nichols are inapplicable to the rigors of pipelaying wherein not only a strong, firm boom structure is required, but also the framework of the sideboom must be sufficiently broad to provide a low enough center of gravity for stability and must include a counterweight to provide sufficient operational stability and safety. It will be readily understood by those skilled in the art that such attributes are not provided by conventional tractors.\nIndicative of initial attempts to improve the sideboom art using hydraulics is an apparatus described by Stefanutti in U.S. Pat. No. 3,265,218. In particular, the use of hydraulically actuated booms and hoist assemblies is described, wherein hydraulic cylinders are used either inside or astride the boom to raise and lower the boom. As will be appreciated by those skilled in the art, the Stefanutti apparatus has not been widely accepted.\nAs taught by Vinton in U.S. Pat. No. 3,938,669, however, while the introduction of such hydraulically-operated sidebooms were anticipated to provide improvements associated with effectively and safely manipulating winches and cables through an operator's interfacing with clutches, brakes, and levers for controlling the position of the boom and for hoisting and lowering a load, such sidebooms failed to provide the prerequisite control and versatility. To attempt to remedy this deficiency in the art, Vinton discloses a hydraulic circuit that includes two separate sources of hydraulic fluid for controlling sideboom movement-related functions. One fluid source provides low volume hydraulic fluid for accomplishing not only precise, low-speed manipulation and control, but also for preventing anti-drift of both boom and hoist. The other fluid source provides high volume hydraulic fluid for providing high-speed operation of these movement and control functions. The plurality of control valves inherent in the Vinton circuit for controlling the hoist and the boom motors are operated via two levers. This apparatus also incorporates a hydraulic cylinder, instead of a drum/cable arrangement for controlling the boom; no drums or cables are involved in raising or lowering either the boom or the hook, except a short cable connected to the hook at the end of a hydraulic cylinder. As is well known by those skilled in the sideboom art, this methodology was a commercial failure.\nAs a further development in the sideboom art, Forsyth teaches in U.S. Pat. No. 5,332,110 a hydraulically-operated sideboom intended to prevent boom over-rotation, to impart positive drive to the boom and load winches, and to provide improved control over free fall and vertical kick-out. These safe operating features are particularly intended for pipelaying applications involving lifting and lowering of large pipes. Indicative of current Caterpillar tractors and sidebooms, embodiments taught by Forsyth are exclusively for newly manufactured hydraulically-operated high-drive tractors which are the antithesis of predecessor low-drive mechanically-operated Caterpillar tractors.\nAs will be appreciated by those skilled in the sideboom art, current Caterpillar hydraulically-operated high-drive sidebooms have a higher center of gravity but operate easier due to less controls than old, discontinued, mechanically-operated, low-drive, Caterpillar sidebooms. Such \"high drive\" models are constructed with an undercarriage that was introduced by Caterpillar in the early 1980's; the sprockets and drive mechanisms for the undercarriage are situated much higher off the ground than the previous conventional \"low drive\" system that was in effect from the inception of the original Caterpillar sideboom. More particularly, the Caterpillar low-drive sidebooms were discontinued from 1972 through 1986 as follows: model 572E was discontinued in 1972; model 572F was discontinued in 1975; model 572G was discontinued in 1986; model 583H was discontinued in 1974; model 583K was discontinued in 1986; model 594G was discontinued in 1975; and model 594H was discontinued in 1986. In addition, Caterpillar hydraulically-operated high-drive sidebooms are very expensive, and some models not only have inherent counterweight obstruction problems, but also are difficult to move from job-site to job-site. While the Forsyth disclosure teaches that his hydraulic pipelayer is adapted for mounting upon a conventional track-laying tractor (i.e., bulldozer), it is well known in the art that the main frame of a sideboom is constructed differently from that of a conventional tractor. In particular, unlike a conventional tractor which is constructed with an oscillating frame, a sideboom is constructed with a rigid frame of wider track gauge than a conventional tractor. Indeed, Caterpillar identifies such tractor and sideboom frames with different serial numbers series. Thus, to obtain the prerequisite performance demanded in the pipelaying art, a drawworks assembly must be mounted upon a frame capable of rigidity to accommodate the pivoting action of a sideboom typically positioned upon rough terrain, with the frame having a sufficiently wide track gauge for stability purposes.\nAs will be understood by practitioners in the art, a drawworks system built upon an old, discontinued, mechanically-operated, low-drive, Caterpillar sideboom--having significantly more controls than a conventional discontinued tractor--inherently suffers from a panoply of problems associated with the simultaneous use of a daunting ensemble of gear-shifting mechanisms, clutches, and brakes, all operated by 6 different hand controls to properly lift and manipulate heavy pipes under conditions generally characterized by unpredictable and adverse terrain. Pipeline construction companies constitute 95% of the users of this type of machinery. As will be appreciated by those conversant with the art, historically, such pipeline construction companies have had to choose between the newer high-drive sidebooms with herein before mentioned faults and high price or the older, discontinued Caterpillar sidebooms that are more economical but are more dangerous and are very difficult to operate. It should also be noted that, as the pipelaying industry continues to mature, the number of skilled sideboom operators has gradually diminished.\nAs should be evident to those skilled in the art, it would be advantageous for construction companies who are continuing to utilize older discontinued Caterpillar sidebooms to have the additional benefits of improved handling, safety, and efficiency. It would be also be advantageous for pipeline contractors to have the ability to expeditiously train low-drive Caterpillar sideboom operators and to simultaneously achieve a level of safety heretofore unmatched by any other sideboom system known in the art, regardless of design. Of course, it would be advantageous for pipelaying contractors to have the benefit of a Caterpillar sideboom that inherently avoids or mitigates the complex levers and the like associated with maneuvering a sideboom, and controlling the lifting and lowering of a pipe load.\nAccordingly, these limitations and disadvantages of the prior art are overcome with the present invention, and improved means and techniques are provided which are useful for effectively and reliably utilizing old, discontinued, mechanically-operated, low-drive, Caterpillar sidebooms for pipelaying applications."}
{"text": "The present invention relates to telephone conference calls, and deals more particularly with conferencing an already-connected telephone call in a seamless manner that is transparent to the call participants.\nWhen two people are participating in a telephone call, they may decide that it would be useful if one or more other persons were also participating. For example, the current call participants may be work colleagues who are discussing an urgent problem, and they might decide that one or more of their co-workers could provide valuable insight into the problem. Currently, the existing telephone call has to be terminated, and a new telephone call has to be initiated using a conference bridge."}
{"text": "The demand for e-commerce brings a unique set of challenges to network infrastructures. For example, one server may be inadequate to provide the required capacity and scalability to serve the increasing size of e-commerce transactions. Server load balancing (SLB) was developed to overcome this problem, where a number of servers (server farm) act as a single server and a special device (dispatcher) dispatches requests to servers in a manner that balances the load on all the servers in the server farm. Some SLB schemes may dispatch requests to the server having the least load, and other SLB schemes may dispatch requests to any server such that the network is optimized, for example. For purposes of simplicity, the server selected by any given SLB scheme is referred to herein as the “best server”.\nIn e-commerce transactions hosted on a website, however, SLB introduces a new problem. When a customer sends information to a website that utilizes an SLB scheme, the information is sent to the best server computed by the dispatcher, Server A, and stored locally on that server. For example, a customer of Amazon.com® selects books and places them in the shopping cart, and then decides to visit other areas of the site before purchasing. When the customer returns to the shopping cart, effectively to access state information (the state of previously stored information), the SLB scheme again computes the best server, this time Server B, to direct the customer's request to. Since the distribution of requests amongst the servers in the server farm may have changed since the customer's last request, the best server on the subsequent request is different from the best server on the first request. Consequently, the customer's information does not exist on Server B, and the shopping cart may be empty when the customer returns to it.\nOne current solution to this problem is to globally maintain state information, such that the state information can be accessed by any server in the server farm. For example, state information can be maintained in a special state server, or even one of the servers in the server farm. One of the disadvantages to this, however, is the latency associated with memory accesses by the servers, as well as dispatcher latency associated with storing state information."}
{"text": "1. Field of the Invention\nThe present invention relates to apparatus and method for verifying an integrated circuit pattern formed on a substrate.\n2. Description of the Related Art\nIn the manufacture of semiconductor devices, there is a need of verifying whether an integrated circuit pattern formed on a substrate is formed as designed in individual processes.\nFIG. 15 is a view to explain a method of verifying a hole pattern. According to the method, a target data TD of preset hole pattern is overlapped with a contour data CD extracted from hole pattern formed on a substrate. In the overlapped state, distances from several places of the foregoing two patterns are measured. It is determined whether or not the difference between the measured distances is within a predetermined allowable value, and thereby, it is determined whether the formed pattern is non-defective or defective.\nThe method of evaluating a pattern shape has been proposed (for example, JPN. PAT. APPLN. KOKAI Publication No. 2002-31525). According to the method, the distance between CAD (Computer Aided Design) output data and pattern data measured by SEM (Scanning Electron Microscope) is measured in the same manner as described above. Based on the measured distance, evaluation is made with respect to the pattern shape.\nThe method of evaluating a pattern defect information has been proposed (for example, JPN. PAT. APPLN. KOKAI Publication No. 2002-93875). According to the method, the error dimension between CAD output data and shape simulation output data is measured. If the measured error dimension is not within a predetermined value, pattern defect information is displayed.\nIn addition, the following method has been proposed (for example, JPN. PAT. APPLN. KOKAI Publication No. 2000-182912). According to the method, comparison is made between an inspection reference pattern formed based on the design layout pattern and the finished predictive pattern. By doing so, pattern distortion of the finished predictive pattern is detected, and thereafter, the detected pattern distortion is identified in accordance with criticality.\nAccording to the method of measuring the distance between the foregoing two patterns, a great many of places is measured, and thereby, the measurement accuracy is improved. However, if the measured places are too many, long time is spent for measuring these places. According to the method, the difference from the corner of the pattern is not so taken into consideration. For this reason, when verifying a pattern having a complicated shape, it is very difficult to determine whether the pattern is non-defective or defective. Moreover, there exist various patterns requiring to bridge neighboring patterns depending on pattern, or to control only line width of the shorter side of the pattern. In the foregoing patterns, the allowable (variable) range for the target dimension has not been clarified. For this reason, according to the conventional method, it is determined that patterns formed with accuracy having no problem on the device operation are defective.\nAccording to the conventional method, time is taken to verify the pattern, and the features of the pattern are not considered; as a result, the pattern is not properly verified. In addition, according to the conventional method, it is difficult to evaluate a complicated pattern."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for taking out laundry articles from a laundry pile, wherein a minor batch of laundry articles is seized by a pull gripper intended therefor and pulled out of the laundry pile and lifted clear above the same.\n2. Description of the Related Art\nSuch machines are known today for use in e.g. industrial laundries where the laundry articles, such as bedclothes and cloths which are, for instance following washing and tumble drying, in a more or less entangled pile, in a damp state.\nThese laundry articles are subsequently to be introduced individually into machines intended e.g. for straightening, ironing and folding the laundry articles.\nIn the laundries, it is thus a task to separate the laundry articles from the pile in which they are following washing and tumble drying in order to allow them to be introduced in to the subsequent machines for e.g. straightening, ironing and folding.\nToday this process is carried out more or less manually but efforts are made to increasingly automate the process, the work positions involved in the manual process being quite straining, and therefore machines have already been developed to carry out at least some steps of the process automatically.\nSuch machines are known from e.g. DE 43 30 911 which features an arrangement consisting of two separate grippers which may cooperate to take out an individual laundry article from the laundry pile and thus to convey the individual laundry article to the subsequent machinery. The grippers in the known machine are so arranged that the gripper which pulls the individual laundry article out of the laundry pile has been optimised accordingly, whereby the gripper seizes a. comparatively small corner or a comparatively small portion of the laundry article with a view to ensuring that the gripper takes out only one laundry article from the pile.\nIt is a problem, however, with these prior art machines that the laundry articles in the laundry pile are often very entangled as a consequence of the washing and drying processes, and that this is enough to require a powerful pull in the laundry articles in order to separate them from the pile, Which is further aggravated by the fact that frequently the laundry articles are slightly damp whereby the friction between the laundry articles is often increased. This may mean that the gripper which takes out the laundry article from the pile is incapable of performing this task in a satisfactory manner since, advantageously, the seizing force and seizing area of the gripper is as small as possible as it is desired to ensure to the highest degree possible that the laundry articles are taken out of the pile individually. If the seizing force is increased, tearing of the laundry articles may as it is desirable to have a small seizing area.\nIn practice this may mean that the gripper cannot take out any laundry articles whatsoever which would require interference from an operator or readjustment of the machine.\nU.S. Pat. No. 4,943,198 features an apparatus for automatically taking out laundry articles from a laundry pile whereby automation of this step of the process is allowed, and whereby it is ensured to a higher degree that no operational stoppages or tearing of the laundry articles occur.\nThis is obtained by a smaller batch of laundry articles being pulled out of the laundry pile by a gripper intended therefor, following which this minor batch of laundry articles is taken away from the laundry pile and conveyed above a support, following which the laundry articles are deposited on the support. This minor batch of laundry articles will then be significantly less entangled than it was in the pile following e.g. washing and tumbling, thereby facilitating the subsequent taking out of individual laundry articles from this minor batch of laundry articles substantially.\nThe pull gripper may thereby be so designed that it seizes a large portion of the laundry articles and in a powerful manner, since there is no requirement that the gripper takes out one single laundry article at a time from the pile. When more laundry articles are taken out at a time it follows that the risk of tearing the individual laundry article is also reduced.\nSince the method and apparatus according to U.S. Pat. No. 4,943,198 deal with the desire to ensure that only one single laundry article is taken out at a time, more room is left for providing such convenient pull gripper as to ensure that laundry articles are always taken out from the pile with no regard to the quantity."}
{"text": "It is a problem both to safeguard data that is stored on a computer system and to restore all or portions of this data that are lost or corrupted. Many computer systems have no protection systems in place and the loss of data from these computer systems is irrevocable. Other computer systems make use of attached data backup systems to store a copy of the data that is stored in the computer memory and updates thereto for eventual retrieval to restore data that is lost from or corrupted in the computer system memory. However, the use of these existing data backup systems is laborious and can be confusing to the casual user. Furthermore, in the instance of network connected computer systems, the use of a common shared data backup system wastes a significant amount of memory by storing multiple copies of the same data, since many programs and customer files on the plurality of computers are duplicates of each other.\nIn information technology, backup refers to making copies of data so that these additional copies may be used to restore the original after a data loss event. These additional copies are typically called “backups.” Backups are useful primarily for two purposes. The first is to restore a computer to an operational state following a disaster (called disaster recovery). The second is to restore one or more files after they have been accidentally deleted or corrupted. Backups are typically that last line of defense against data loss, and consequently the least granular and the least convenient to use.\nSince a data backup system contains at least one copy of all data worth saving, the data storage requirements are considerable, which data storage requirements can be exacerbated by the method used to perform the data backup where change tracking is wasteful of memory. Organizing this storage space and managing the backup process is a complicated undertaking. A data repository model can be used to provide structure to the data storage device for the management of the data that is backed up. In the modern era of computing there are many different types of data storage devices that are useful for making backups. There are also many different ways in which these data backup devices can be arranged to provide geographic redundancy, data security, and portability.\nBefore data is ever sent to its data backup storage location, it is selected, extracted, and manipulated. Many different techniques have been developed to optimize the backup procedure. These include optimizations for dealing with open files and live data sources as well as compression, encryption, and de-duplication, among others. Many organizations and individuals require that they have some confidence that the backup process is working as expected and work to define measurements and validation techniques to confirm the integrity of the backup process. It is also important to recognize the limitations and human factors involved in any backup scheme.\nDue to a considerable overlap in technology, backups and data backup systems are frequently confused with archives and fault-tolerant systems. Backups differ from archives in the sense that archives are the primary copy of data and backups are a secondary copy of data. Data backup systems differ from fault-tolerant systems in the sense that data backup systems assume that a fault will cause a data loss event and fault-tolerant systems assume a fault will not cause a data loss event.\nData Repository Models\nAny backup strategy starts with the concept of a data repository. The backup data needs to be stored somehow and probably should be organized to a degree. It can be as simple as a manual process which uses a sheet of paper with a list of all backup tapes and the dates they were written or a more sophisticated automated setup with a computerized index, catalog, or relational database. Different repository models have different advantages. This is closely related to choosing a backup rotation scheme. The following paragraphs summarize the various existing backup models presently in use.\nUnstructured\nAn unstructured repository may simply be a writable media consisting of, for example, a stack of floppy disks or CD-R media with minimal information about what data from the computer system was backed up onto this writeable media and when the backup(s) occurred. This is the easiest backup method to implement, but probably the least likely to achieve a high level of recoverability due to the dearth of indexing information that is associated with the data that is backed up.\nFull+Incremental\nA Full+Incremental data backup model aims to make storing several copies of the source data more feasible. At first, a full backup of all files from the computer system is taken. After that full backup is completed, an incremental backup of only the files that have changed since the previous full or incremental backup is taken. Restoring the whole computer system to a certain point in time requires locating not only the full backup taken previous to that certain point in time but also all the incremental backups taken between that full backup and the particular point in time to which the system is supposed to be restored. The full backup version of the data is then processed, using the set of incremental changes, to create a present view of the data as of that designated certain point in time. This data backup model offers a high level of security that selected data can be restored to its present state and this data backup model can be used with removable media such as tapes and optical disks. The downside of this data backup process is dealing with a long series of incremental changes and the high storage requirements entailed in this data backup process, since a copy of every changed file in each incremental backup is stored in memory.\nFull+Differential\nA Full+Differential data backup model differs from a Full+Incremental data backup model in that after the full backup is taken of all files on the computer system, each incremental backup of the files captures all files created or changed since the full backup, even though some may have been included in a previous partial backup. The advantage of this data backup model is that restoring the whole computer system to a certain point in time involves recovering only the last full backup and then overlaying it with the last differential backup.\nMirror+Reverse Incremental\nA Mirror+Reverse Incremental data backup model is similar to a Full+Incremental data backup model. The difference is that instead of an aging full data backup, followed by a series of incremental data backups; this model offers a mirror that reflects the state of the computer system as of the last data backup and a history of reverse incremental data backups. One benefit of this data backup method is that it only requires an initial full data backup. Each incremental data backup is immediately applied to the mirror and the files they replace are moved to a reverse incremental backup. This data backup model is not suited to the use of removable media since every data backup must be done in comparison to the data backup mirror version of the data. This process, when used to restore the whole computer system to a certain point in time, is also intensive in its use of memory.\nContinuous Data Protection\nThis data backup model takes the data backup process a step further and instead of scheduling periodic data backups, the data backup system immediately logs every change made on the computer system. This is generally done by saving byte or block-level differences rather than file-level differences. It differs from simple disk mirroring in that it enables a roll-back of the log and thus can restore an old image of data. Restoring the whole computer system to a certain point in time using this method requires that the original version of the data must be processed to incorporate every change recorded in each differential change to recreate the present version of the data.\nProblems\nIn spite of all of these various methods of data backup, existing data backup systems (including both hardware and software) fail to ensure that the user can simply plug in to the computer system to “back-up” the data stored therein, and also enable recovery of a revision of a file from a point-in-time, and enable all of the hard disk(s) in the computer system to be restored to a point-in-time. Existing data backup systems fail to efficiently track and store the state of multiple file systems over time, while allowing for correct disk-level and file-level restoration, to a point-in-time, without storing a significant amount of redundant data. These data backup systems require the user to learn new technology, understand the file system of the computer system, learn how to schedule data backup sessions, and learn new controls that must be used for this new functionality. Furthermore, the restoration of lost files is difficult using these data backup systems."}
{"text": "Computer control of a machining system that involves the CAD/CAM based support has been widely accepted to improve productivity and reduce production cost. Recently, more intelligent functions have been developed and integrated into CNC machine tools. CAD/CAM provides the facilities to create and monitor tool paths to use on the workpiece. In some CAM software programs, machine tools and machine virtual environments can be utilized to dynamically simulate the machining operations. These dynamic simulations provide NC program generation and verification, material removal analysis and collision detection error. With the process simulation, the tool path can be analyzed and verified before actually machining the part. It has become easier to machine complex parts more accurately and more quickly with the advancement of simulation tools. However, in selection of machining strategies, the methods offered by CAM software often are based on the parts' geometrical information with little or no consideration of the machine tool capability or the physics of metal cutting. On the other hand, machine tools (or operators) have limited information about NC programs, hence it is difficult to judge whether machining is performed properly. Running machine tools under undesirable operating conditions can cause damage to tools, the machine tools or workpiece. Operating at, near or over the machine limits, for a short time or over a long period of time, can lead to damage to the tool, the machine tool or workpiece.\nMoreover, in practice, operating parameters are still mainly selected based on either on machining handbooks and/or tool manufacturer's catalogues which are typically very conservative and aggressive, respectively. Therefore, it has been difficult to perform the machining under an optimal condition, which either leads to low productivity or deterioration in machining accuracy and surface roughness. Moreover, when the tool is a cutting tool which is engaged in cutting with a rapidly increasing cutting load, damage easily occurs to the cutting tool as well as work material to be machined."}
{"text": "Commuters in big cities in the United States, e.g., Washington, District of Columbia, Los Angeles, San Francisco, New York City, Boston, Houston, Atlanta, Chicago, Philadelphia, Seattle, Miami, Dallas, Detroit, San Diego, Phoenix, etc., spend an average of 50 hours a year waiting in traffic. In one study, the cost of all this wasted time and fuel is estimated to be more than $100 billion a year. In other big cities around the world, e.g., Beijing, Shanghai, Tokyo, Brussels, Antwerp, Milan, London, Paris, Rotterdam, etc., the average can be much higher. A context-based traffic flow control system and method thereof can advantageously help minimize travel time and decrease traffic congestion."}
{"text": "(1) Field of the Invention\nThe present invention relates to a roller assembly for a roller furnace for heat treating slabs of steel, etc. The roller assembly has a number of spaced carrier rings for carrying the slabs and thermal insulation between the carrier rings. The slabs are transported through the roller furnace by the roller assemblies on the carrier rings for heat treatment thereof. The hollow cavity of the roller assembly is usually water cooled. The thermal insulation is for protecting the roller assembly from the hot slab, as well as preventing the hot slab from cooling due to the presence of the roller assembly.\n(2) State of the Prior Art\nDE 32 31 736 C2 discloses a sheathing for insulating a cooled slide, skid or transverse pipe in a furnace. Moldings for the bottom side of the pipe made of a ceramic fiber material and moldings for the upper side of the pipe made of a refractory ceramic are provided for adapting the bottom and upper sides for the different stresses the sides will undergo. This type of an arrangement is unsuitable for a rotating roller assembly. Preformed insulating moldings result in joints in the thermal insulation of the roller assembly. Such joints have a negative impact on the insulating effect of the thermal insulation.\nMoldings for insulating pipes are also described in DE 31 25 440 A1 and in DE 36 09 047 A1. The types of moldings described in these disclosures are unsuitable for roller assemblies, since the moldings have very little ability to withstand the dynamic stresses of a rotating roller assembly, and since the joints in the moldings negatively impact on the effect of the thermal insulation."}
{"text": "In the manufacturing steps of rubber products such as tires, a strip-shaped belt member is used that is made of a plurality of steel cords disposed side by side at a predetermined cord angle and covered with rubber. This belt member is formed by cutting the member obtained by covering the steel cords disposed side by side with rubber, at a predetermined length and at a predetermined angle set in accordance with the cord angle, and bonding together a plurality of the cut members (See Japanese Unexamined Patent Application Publication No. 2009-226817A, for example). This strip-shaped belt member is wound around a molding drum and molded into a cylindrical shape. Further, a tread rubber forming a tire tread is molded into a cylindrical shape by winding strip-shaped unvulcanized rubber around the molding drum.\nWhen the strip-shaped belt members and the strip-shaped unvulcanized rubber are wound around the molding drum in a meandering manner, a cylindrical member having strain in the width direction of the drum is molded. Further, in the cylindrical member molded by winding the plurality of strip-shaped members around the molding drum, also when the strip-shaped members are offset with respect to each other while being misaligned in the width direction by an amount greater than a set value, the strain in the width direction of the drum occurs in the cylindrical member. Manufacturing tires using the cylindrical member subject to this kind of strain makes it difficult to ensure consistent tire quality. Therefore, it is important to accurately assess a widthwise misalignment of the cylindrical member obtained by molding the strip-shaped member in the cylindrical shape."}
{"text": "1. Field of the Invention\nThis invention is concerned with preparing alginic acid esters known as \"alkylene glycol alginates\" by a procedure for treating seaweeds of the Class Phaeophyceae. This novel and improved process involves treatment of the seaweed with an acid and thereby forming alginic acid therein and subsequently reacting the alginic acid content of the treated seaweed, with an alkylene oxide to form an alkylene glycol alginate in situ without prior isolation of alginic acid or its salts from the solid seaweed residue.\n2. Description of the Prior Art\nAlginic acid is a polyuronic acid generally believed to consist of two uronic acids; mannuronic acid and guluronic acid, the proportions of which vary depending on factors such as, for example, seaweed species, plant age and seasonal variation. Alginic acid in the form of mixed water insoluble salts, in which the principal cation is calcium, is present in seaweeds of the Class Phaeophyceae, typical examples of which are: Fucus vesiculosus, F. spiralis, Ascophyllum nodosum, Macrocystis pyrifera, Alaria esculenta, Laminaria longicruris, L. digitata, L. saccharina, and L. cloustoni.\nMethods for the recovery of water insoluble alginic acid and its water soluble salts, particularly sodium alginate, are well known. The first extraction process was patented by Stanford in British Pat. No. 142 (1881). A series of variations of Stanford's method have subsequently been described in the patent literature. The most recent and familiar are the processes of Green, U.S. Pat. No. 2,036,934, and Le Gloahec and Herter, U.S. Pat. No. 2,128,551.\nIn Green's process, seaweeds of the Class Phaeophyceae are treated with dilute acid, such as dilute hydrochloric acid, followed by water washing. This converts the natural alginate salts present to alginic acid. Salts, residual hydrochloric acid, and water soluble organic materials are then removed by washing. The pretreated seaweed is then chopped or milled and an excess of sodium carbonate added together with a quantity of water to extract the water soluble sodium alginate formed. The mixture is filtered to recover the clarified sodium alginate solution to which a solution of calcium chloride is added to precipitate the alginate in the form of water insoluble calcium alginate. The highly hydrated precipitate is separated from the solution of sodium chloride, excess calcium chloride, and soluble color bodies, particularly phenolic compounds, usually present in the original extract. The precipitate may still retain some color bodies, even after water washing, and calcium hypochlorite is added as a bleaching agent. The precipitate is then treated with dilute hydrochloric acid to convert the purified calcium alginate to alginic acid. This gelatinous precipitate is washed with water to remove excess hydrochloric acid and calcium salts. It may be neutralized to provide the purified sodium alginate of commerce.\nIn the Le Gloahec and Herter process, seaweeds of the Class Phaeophyceae are treated with a solution of calcium chloride and the water soluble components of the seaweed are then extracted and removed by draining. The seaweed is then treated with dilute hydrochloric acid, drained, and water washed. Sodium carbonate is added and the seaweed mixture milled and diluted with water to extract the sodium alginate. The slurry is aerated to separate insoluble materials by a method of air flotation and the residual clarified solution is treated with a decolorizing agent. The purified sodium alginate could then be recovered by conventional means, as for example, by alcohol precipitation. To recover alginic acid, the solution is treated with dilute sulfuric acid, the precipitate being separated and pressed to remove the aqueous solution containing sodium sulfate and excess sulfuric acid. The precipitate is dehydrated with alcohol, washed with further quantities of alcohol, and then dried to produce the alginic acid of commerce.\nMethods for preparing alkylene glycol alginates are also described in the patent literature. Steiner in U.S. Pat. No. 2,426,125 discloses a method for the manufacture of glycol alginates followed subsequently by Steiner et al in U.S. Pat. Nos. 2,494,911 and 2,494,912, Nielsen et al in Canadian Pat. No. 904,847, and Pettitt et al. in their related disclosures in Canadian Pat. No. 942,744 and U.S. Pat. No. 3,772,266. All of these patents describe, with variations in procedure, methods by which alkylene oxides are reacted with alginic acid.\nThe prior art teaches the preparation of glycol alginates by the reaction of alkylene oxides with the alginic acid produced by methods known in the art, such as those described above or those mentioned by Steiner in U.S. Pat. No. 2,426,125, viz. U.S. Pat. Nos. 1,814,981 to Thornley et al., 2,036,922 to Clark et al. and 2,036,934 to Green. It will be noted that Steiner U.S. Pat. No. 2,426,125 states (column 3, lines 3-6) \"that the free acid from the above or other modifications of the general method may be used for my purpose provided only that the acid be in a state of commercial purity\". The wet alginic acid is typically partially dried so as to contain approximately 50% by weight of water. This acid is milled to provide a large surface area; and either before, during, or after the drying process, it is partially neutralized with any suitable base so that between 5 and 30% of the carboxyl groups of the acid are combined with the base. It is essential that the base used be thoroughly disseminated through the acid so as to avoid the possibility of part of the acid being completely neutralized, and thus rendered unsuitable for the esterification reaction, while some of the acid remains free of partial neutralization. For this reason, Steiner describes several methods of partial neutralization. For example the base may be dispersed in a quantity of low boiling alcohol or ketone and added as a slurry to the stirred, wet, milled alginic acid. A wetting agent may also be added to aid dispersion.\nSteiner's partially neutralized, finely divided, alginic acid is reacted with an alkylene oxide such as ethylene oxide or propylene oxide (1,2-epoxypropane). In commercial practice, the preferred glycol alginate is propylene glycol alginate formed by reaction of alginic acid with propylene oxide. The reaction proceeds with both the formation of the glycol alginate, and hydrolysis of the alkylene oxide due to the presence of water and the low pH of the mixture. Thus, a greater quantity of the alkylene oxide than the stoichiometric equivalent is required. It is believed that molar ratios of between 2:1 and 3:1 alkylene oxide: alginic acid are used, although Nielsen et al in Canadian Pat. No. 904,847 state (page 9, lines 17-18) that \"the amount of propylene oxide may be from 1 mol to about 25 moles per mol alginic acid.\" In Nielsen's case, the reaction is carried out in a diluent, such as an alcohol or ketone, and a substantial portion of the oxide is thus not consumed in the esterification reaction. It is generally recognized that complete reaction of oxide and available carboxyl groups is not only difficult, requiring the use of large quantities of alkylene oxide, but unnecessary.\nSaid Pettitt et al patents disclose partly neutralized alginic acid of critical 65-78% solids content reacting with propylene oxide gas in the absence of air in 3 hours or less at 60.degree.-100.degree.C. In the examples, the oxide : acid molar ratios range from 2.8:1 to 15:1.\nWallerstein et al. in U.S. Pat. No. 2,478,988 teach the use of propylene glycol alginate as a foam stabilizing agent. According to Steiner U.S. Pat. No. 2,659,675, the glycol alginate used for the purpose of the Wallerstein patent was designed for use as an emulsifying agent in French dressings and the like and had the analysis: neutralization as ammonium alignate 30 to 40%; esterification 25 to 40%; and unreacted acidity 20 to 45%. In the same patent, Steiner teaches that an improved foam stabilizing agent has the following characteristics among others: neutralization as sodium alginate 15 to 20%; esterification 65 to 80%; and unreacted acidity 5 to 15%. Therefore, depending on the end use of the glycol alginate product, it appears that a fully satisfactory product may have any of the following characteristics: neutralization of between about 5 to 40% with any suitable base or mixture of bases; esterification of from about 25 to 80%; and unreacted acidity of from about 5 to 45%. In general the glycol alginate should be completely soluble in water giving a solution the pH of which is not less than about 3.5. Storage stability of the dry product may be impaired if the pH of its aqueous solution is below 3.5.\nFrom the above description of the prior art, it is quite evident that prior art techniques for the preparation of glycol alginates have involved formation of alginic acid, separation thereof from seaweed residue and subsequent esterification of the alginic acid to produce the glycol alginate and have been rather complex, slow and costly in terms of chemical and processing requirements. Of the procedures described, the initial preparation of alginic acid appears the most complex, although separating the alginic acid in a form considered to be suitable for esterification, as described earlier, also appears to require a great deal of care and attention.\nIt has now been discovered that, contrary to previous belief, the alginic acid used for reactions, such as the esterification reaction described above, need not be in the extracted state of commercial purity. In other words, the present invention is based on the discovery that it is not necessary to isolate the alginic acid from the solid residue of the precursor seaweed in order to react the alginic acid with an alkylene oxide. In addition, the glycol alginate products of my novel process are capable of meeting the various commercial requirements and specifications such as those outlined above, together with color, purity, etc."}
{"text": "The invention relates to a front, side or rear view monitoring method for a motor vehicle (such as an automobile) which can warn a driver by using images taken by a video camera installed on the front, side or rear part of the motor vehicle and sensing a motor vehicle or an obstacle ahead of the driver's own motor vehicle during a drive.\nA conventional method of this type is disclosed in a publication of Japanese Unexamined Patent Publication No. 241855/1990. The method disclosed in the publication comprises the steps of: picturing a front view from a driver's own motor vehicle that is moving; recognizing a movement of a single point in the pictured front view as an optical flow at predetermined times; gathering information, i.e., the position of a preceding motor vehicle relative to the driver's own motor vehicle and the relative speed, based on the optical flow as well as on the distance of the preceding motor vehicle as measured by a range sensor; and warning the driver of danger when the presence of danger is judged.\nTo calculate the optical flow, a method called a matching method has here ore been adopted. The matching method detects a corresponding point in two images. The matching method comprises the steps of: setting a window W1 for a target pixel P in an image at a time t, as shown in FIG. 10 (a); calculating a correlative value while moving the window over the entire part of the image or a neighboring area; and finding a corresponding point in a window W2 defined at the time the correlative value is maximized, i.e., a pixel Q, as shown in FIG. 10 (b). An arrow PQ indicates an optical flow. To calculate the above-mentioned correlative value, the following equation is used. EQU .SIGMA.(W1.sub.(x, y) x W2.sub.(x, y))/(.SIGMA.W1.sup.2.sub.(x, y) x .SIGMA.W2.sup.2.sub.(x,y)).sup.1/2 ( 1)\nwhere W1.sub.(x, y), W2.sub.(x, y) are the outputs of coordinates (x, y) within the windows W1, W2.\nHowever, in order to judge danger, the conventional method requires a range sensor for measuring the distance of a preceding motor vehicle to identify the position of the preceding motor vehicle relative to the driver's own motor vehicle, in addition to a video camera for picturing the front view.\nFurther, in order to calculate the optical flow, the operation of detecting the corresponding point in an image must be performed. As described above, such operation comprises detection of all the pixels in the entire part or a neighboring area of the image relative to a single pixel within the image to find a single corresponding point. To detect the corresponding point over the entire part of the image, this operation must be performed on all the pixels. In addition, a great number of calculations are required to be performed to find correlative values that are to be used as indexes for search.\nTherefore, an increased volume of calculation makes realtime processing impossible. For achievement of the realtime processing, a high-speed processor is required, which elevates hardware costs. Further, it is not easy to distinguish the source of an optical flow; e.g., whether the optical flow is derived from the preceding motor vehicle or from scenes outside the road, or characters, symbols, or white lines on the road surface, or the like."}
{"text": "The concept of the “intelligent network” (IN) was developed for efficient deployment of resources when using digital communication networks. This is a structural concept for flexible service control and service development through the use of software and functionality implemented therewith. Internationally IN is defined in the Q.1200 series recommendations, for which the ITU (International Telecommunications Union) is responsible. The basic thinking behind IN is the separation of switching functions and service functions, i.e. resources for connection set-up are no longer located wholly in the individual switching point but are centralized. The advantage of this concept is simpler management and adaptation of these centralized resources. For example, when new services are introduced, not all switching points have to be adapted individually.\nThere are three hierarchical levels in the IN architecture, management level, service control level, and switching and transport level. The switching points—generally referred to as SSP (service switching point)—of the lowest level, i.e. the switching and transport level, are controlled by control points—generally referred to as SCP (service control point)—at service control level for the implementation of services and where appropriate are supplied with data from centralized databases. SSP and SCP are defined as physical entities for the IN concept in Q.1205.\nGeneral aspects of an INAP (intelligent network application protocol) protocol are described in Q.1208 for IN applications. At present INAP protocols with CS2 (capability set 2) and CS3 (capability set 3) operational features in accordance with the ITU recommendations numbered Q.122x and Q.123x are frequently used for IN applications.\nFurther specifications are required to provide telecommunication services and service features in the context of the IN concept. For example, for ISDN (integrated services digital network) applications signaling functions are standardized in the context of the ISUP (ISDN user part) in ITU documents Q.761-Q.764 and Q.766, with the ISUP representing a protocol based on signaling system no. 7 (SS#7), which allows separation of the transfer of user data and control data and which is standardized by the ITU in the Q.700 series. For ISDN applications in the context of the IN concept the ITU has produced standards Q.1600 and Q.1601 which govern the interaction of ISUP and INAP CS1 or INAP CS2.\nGiven the increasing importance of the transfer of user data and any signaling data via different networks or carriers, as is frequently the case, for example, with IP based applications, signaling protocols are being developed, which allow a specific set of services and service features, regardless of the carrier. An important example of such a service is internet telephony (VoIP), in which connection-based voice transfer is to be achieved via different physical networks. The ITU draft of Q.1901 describes a signaling protocol, the BICC (bearer independent call control) protocol, with which the services and service features defined for the ISDN network—hereafter referred to as ISDN services or service features—are to be implemented, regardless of the carrier. The BICC protocol is based on the ISUP protocol. As described in the draft of Q.765.5, the transport mechanism (application transport mechanism) set out in Q.765 can be used for the transport of carrier-related data using the BICC protocol.\nWhen implementing ISDN services and service features in the context of the IN concept with user data transfer via a packet network using a carrier-independent signaling protocol, such as, for example, BICC CS2, there is the problem that the parameter changes necessary or desirable for this are not supported."}
{"text": "This invention relates to the spraying art and, more particularly, to a spray manifold in and through which the quantity and rate of flow (or, more accurately, the spray) of fluid is easily adjustable by moving a pintle.\nSpray manifolds and pintles are well known. Equally well known are the facts: that the desired flow of a fluid into and through a spray manifold is conventionally adjusted by deforming the manifold; and, that any exposure of the deformed manifold to any hot operating environment adversely alters the deformation and, thus, the already-adjusted flow through the manifold, thereby requiring further adjustment of the manifold.\nWe have eliminated this re-adjustment problem with a uniquely stuctured spray manifold, wherein the desired flow into and through the manifold is very easily adjusted and maintained by moving a pintle, instead of further deforming the manifold, as is done traditionally.\nThereby, we have significantly advanced the state-of-the-art."}
{"text": "1. Field of the Invention\nThis invention relates to an addressing circuitry and more particularly to an improved addressing circuitry for driving dot matrix vertical scan display apparatus, such as a plasma display panel.\nOf the number of display apparatus, dot matrix vertical scan plasma display panels have been of a relatively recent introduction. One such manufacture of plasma display panel has been Burroughs Corporation. Burroughs has introduced a number of different models. These displays usually have a number of rows and a number of characters per row to display ASCII characters. In earlier models, number of the rows and characters per row have been based on binary progression numbers; for example, the number of rows of the characters displayed on the panel is one row or two rows or four rows or eight rows which is in the progression of 2.sup.0, 2.sup.1, 2.sup.2, 2.sup.3 and likewise the number of characters per row is 16 or 32 characters or 2.sup.4, 2.sup.5. Each of the characters is represented by illuminating selected ones of seven by seven dot cross points. More recently a new line of models have been introduced wherein the number of rows and the number of characters per row are not binary progression numbers. The new models have been introduced to meet the market needs.\nWhere the display panels wherein the number of rows and the number of characters per row correspond to a binary progression number, the prior art addressing circuitry provided by manufacturers of the displays have been found satisfactory.\nHowever, the present inventors found that this is not the case with the displays having character rows and having characters per row which are not of binary progression numbers. More specifically, the present inventors found that the prior art addressing circuitry if used to operate such a display, would read out the characters in the random access memory but not in sequence. Instead, they skipped certain character locations in the random access memory in the read out process. This wasted random access memory locations.\nThe prior art addressing circuitry is also found not to be completely satisfactory in providing cursor, blanking and blinking operation of the display."}
{"text": "This relates generally to compression of depth buffers.\nPrevious depth buffer compression schemes compress the depth obtained by rasterizing static triangles. This provides generous bandwidth usage savings, and so is of great importance to graphics processors. However, stochastic rasterization for both motion blur and depth of field is becoming a reality even for real-time graphics, and previous depth buffer compression algorithms fail due to the irregularity of the positions and depth of the rendered fragments."}
{"text": "1) Field of the Invention\nThis invention relates to automobile passenger protection, and specifically, is a method and apparatus for modifying and adjusting existing shoulder belts to fit passengers who are currently unable to wear them.\n2) Discussion of Prior Art\nStandard three-point automobile seat belts are only designed to restrain a passenger of average adult size. The shoulder portion crosses the passenger seat, and hence the body of the passenger seated thereon, at a location which is fixed in relation to the passenger seat, not to the passenger. The crossover point is not always properly over the center of the passenger's breastbone, a.k.a. the mesosternum. In fact, this crossover point may even be over the neck or face of a passenger of below average height or size. The result is that such a passenger, e.g., a child, is unable to use standard shoulder belts. It has been thought desirable to adjust the position of the belt and/or passenger so that the shoulder belt would properly cross the mesosternum, and permit the passenger to wear a shoulder belt.\nSome devices of the prior art has attempted to lower the shoulder belt crossover point by lowering the \"vertical transition point\", the point from which the shoulder portion of the belt turns and proceeds diagonally across the passenger. The position of this point is originally located at the \"vertical transition member\", which is the hardware member, typically a D-ring, to which the shoulder portion of the belt rises and meets the vertical portion of the belt. This vertical transition member is typically mounted in the car's interior near the passenger's shoulder. Providing means to move the vertical transition point by moving the actual vertical transition member necessarily required a substantial modification of the vehicle interior. Hence, such prior art devices could not be used on existing automobiles. Furthermore, they are a cumbersome solution to the problem.\nOther prior art attempts attempted to lower the vertical transition point, and hence the point at which the shoulder belt crossed the occupant's chest, without moving the actual vertical transition member. One such approach is shown by Pollitt's U.S. Pat. No. 4,236,755, and Ashworth's U.S. Pat. No. 4,289,352. These provide, along with the standard safety belts, an additional belt mounted about or adjacent to the seat. Further provided is a deflection member, which is adjustably mounted on the belt. The shoulder belt may be deflectively adjusted by the deflection member so as to better cross the occupant's chest. A strong disadvantage shared by Ashworth's and Pollitt's apparatuses is that they each require adding an additional separate, independent belt to the car, upon which their deflection member is mounted. Thus, they cannot use existing belts, but require either retrofitting existing automobiles or redesigning future automobiles to include the extra parts. Furthermore, the deflection members utilized by Ashworth and Pollitt lack a superior means for being securely, yet adjustably, mounted to the belt supporting them, as well as means for optimal torsional positioning of the shoulder belt which they are deflecting so as to lay flat as possible in fullest desired contact with occupant. Furthermore, the devices taught are not substantially co-planar with the interior automobile and belt surfaces, and hence present an additional hazard to the user.\nThus, nothing in the prior art shows, teaches, or suggests the present invention, especially of its: (1) providing a method and means of adjusting the point at which the shoulder belt crosses the passenger by lowering the effective upper end of the diagonal portion of the seat belt by deflecting it towards the vertical portion of the belt itself with a device that is itself frictionally and mechanically bound to the vertical portion of the belt and held to thereby, and (2) providing belt adjustment means which are mounted on the vertical portion of the belt and which engage the shoulder portion of the belt and which are used by sliding them up and down on the vertical portion of the belt so as to raise and lower the effective connection of the shoulder belt to the vertical portion of the belt, and also so as to raise and lower the effective height of the point from which the shoulder belt proceeds diagonally downward and also to raise and lower the point at which the shoulder belt crosses the passenger.\nThe present invention overcomes the disadvantages of the prior art, and has numerous objects and advantages which will become apparent from the drawings and the ensuing description."}
{"text": "1. Field of the Invention\nThe present invention relates to a screwdriver. More particularly, the present invention relates to a screwdriver of the type with a ratchet mechanism allowing a change in the driving rotational direction for driving a screw and in a free rotational direction not driving the screw.\n2. Description of the Related Art\nU.S. Pat. No. 6,047,802 discloses a ratchet driving mechanism for a screwdriver. The ratchet driving mechanism includes a main body having two slots. A detent is received in each slot. The detents are selectively biased outward by springs to engage with teeth of a cylindrical rotary seat, thereby allowing a change in the driving rotational direction for driving a screw and in a free rotational direction not driving the screw. In manufacture, formation of the slots requires two elongate mold cores that are placed in a cavity of a mold. Since each slot does not extend through the main body, each mold core has an end unsupported in the cavity. When molten metal is filled into the cavity, the mold cores sway and wobble due to high temperature and high pressure in the mold such that errors are incurred in the specification of the slots. Namely, if the slots are not processed subsequently, the movement of the detents in the slots would not be smooth, leading to risks of getting stuck and malfunction. The cost is thus increased. Further, only two detents are used, and only one of them is engaged with the teeth of the cylindrical rotary seat such that the detent might slide relative to the teeth or the detent and/or the teeth might be damaged when a large torque is applied to the screwdriver. Hence, such a screwdriver can only be driven by a small torque, leading to limited application of the screwdriver, and if in any event a torque exceeds what the screwdriver can withstand, the detent which moved to a position in engagement with the teeth would be pushed out of the slot by the teeth.\nU.S. Patent Application No. 2006/0096422 discloses a ratcheting mechanism for a driving tool. The ratcheting mechanism is utilized for permitting the driving tool to make a change in the driving rotational direction. The ratcheting mechanism includes two pairs of detents and each pair has a spring connected therebetween. The two pairs of detents are received in two spaces respectively defined in a driving member and each detent has one side abutted and supported by a wall of the space and an opposite side selectively engagable with a plurality of teeth defined in a control ring. Additionally, the control ring is rotatable and operably moves one of one pair of detents to a position in engagement with the teeth and has a stir bar that moves the other of one pair of detents to a position away from the teeth. Thus, when the driving tool is rotated in one driving direction two detents (i.e. upper left and bottom right or upper right and bottom left) would engage with the teeth at the same time, and this results that the driving tool is more capable of withstanding higher driving torque in comparison to the prior admitted invention. Further, while the two detents engage with the teeth, the other two detents are kept away from the teeth.\nBut, it is noticed that this invention has a practical difficulty because in operation when the control ring is rotated counterclockwise, the upper right and bottom left stir bars cause the upper right and bottom left detents to disengage or separate from the wall of the space and coact with the teeth, and even worse, after subsequent direction changes the detents will not mesh with the teeth correctly. Likewise, when the control ring is rotated clockwise, the upper left and bottom right stir bars would cause the upper left and bottom right detents to disengage from the wall of the space and coact with the teeth, leading to incorrect meshing after subsequent direction changes."}
{"text": "This invention relates generally to a computer network-based system and more particularly to systems and methods for managing workflow both internally to provide services to customers and externally to allow customers to monitor workflow.\nElectronic information exchange has improved workflow management problems to some extent, however differences between, for example, two companies interested in making a deal still exist. Legacy formatting issues between the companies and differences in procedures have left much of the inefficiencies found in paper based methods unaffected. For example, time expended manually in gathering information regarding the deal and possibly converting the information into a different format contribute to the inefficiencies. Overall, both the traditional paper methods and known electronic methods can be very frustrating and may not necessarily facilitate resolving workflow management issues.\nIt would be desirable to manage workflow issues by implementing a web-based system that allows real-time access to project status, both internally to the supplier of services and externally to the customer where they can monitor the workflow of services provided, while eliminating the inefficiencies of known systems thereby providing higher customer satisfaction and reduction in deal cycle time.\nA workflow management system providing an integrated approach to exchanging information real-time and notifying responsible persons to pending commitments concerning a particular deal is disclosed. The web-based system includes a server having a centralized database of deal data. The connecting of at least one client system to the server facilitates a method of uploading initial proposed deal data, notifying an underwriter of the uploaded proposal, uploading a workflow timeline for the proposed deal and notifying responsible persons that actions for tasks are due according to the timeline.\nThe workflow management system is further configured to restrict database access only to authorized users for viewing of and providing updates to the deal."}
{"text": "The following U.S. patents are believed to represent the current state of the art:\nU.S. Pat. Nos. 5,496,236; 5,016,870; 4,860,763; 4,550,908 and 4,402,502"}
{"text": "Polishing is an activity which is useful to maintain the appearance and condition of footwear, and other accessories such as belts, purses, and the like. These objects, whether made from fabric, leather, or similar materials, often become dirty or dull as a result of use. Various implements have been developed for applying treatment agents to the surfaces of such objects to clean or enhance their finish, and for polishing after treatment. In many cases, more than one implement is required to complete the polishing process; for example, when polishing shoes, it is often necessary to use one implement for applying polish to the shoes, and another implement to buff the polish in order to achieve the desired shine. Moreover, such implements generally require both hands, and use of them often results in transfer of treatment agent to the hands of the user. Because of the need to use more than one implement, and the mess that is often associated with the process, many users are not inclined to regularly polish their shoes. Accordingly, there is a need for a shoe care apparatus which is multifunctional, can be used with one hand to achieve effective treatment and buffing, and which minimizes the need for the user to come in direct contact with polish and other treatment agents."}
{"text": "There are many circumstances wherein individuals can become overheated. Vigorous exercise is a clear example, as well as leisure activities under conditions of high temperatures and/or intense sunlight. Active and effective cooling of individuals under these circumstances can protect their health and significantly increase their comfort and enjoyment.\nEvaporative cooling is well known as a highly effective means for cooling individuals. Indeed, it is the mechanism by which the body cools itself through perspiration. The discomfort and potential dehydration of cooling by perspiration can be avoided through the application of a mist of water to an individual, which cools the body in essentially the same manner as perspiration, and can be even more effective than perspiration since the mist impacts the skin at a temperature significantly below body temperature, and hence absorbs more heat than an equivalent quantity of perspiration.\nMeans for generating and applying a water mist to one or more individuals are well known. However, misting does not always provide sufficient cooling efficiency, due to limited evaporation rates. Also, misting can lead to wetness if the evaporation rate is not sufficient."}
{"text": "1. Field of the Invention\nThe present invention relates to a blood pressure measuring apparatus including a surface-pressure sensor adapted to be pressed against an artery of a living body or subject via a body surface of the subject.\n2. Related Art Statement\nThere is known a blood pressure (BP) measuring device which non-invasively measures a BP value of a living subject based on a pulse-sound signal or a pulse-wave signal which is continuously obtained while an inflatable cuff wound around subject's body portion such as an upper arm is inflated, or is deflated after being inflated, to press the body portion. The non-invasively measured BP value is an approximation of an intra-arterial BP value directly or invasively measured from inside an artery of the subject.\nIn the conventional BP measuring device, however, the cuff presses the subject's body portion all around with a considerably great force, so that the subject may suffer the discomfort due to the use of the cuff.\nMeanwhile, Unexamined Japanese Utility Model Application laid open under Publication No. 1(1989)-161707 discloses a blood pressure (BP) monitor including (a) an inflatable cuff, (b) a fluid supply which supplies a pressurized fluid to the cuff, (c) a control device which changes the fluid pressure in the cuff, (d) a surface-pressure sensor adapted to be pressed on subject's body surface to detect a pressure pulse wave produced from subject's artery and generate a pressure signal representing the detected pulse wave, and (e) BP determining means for determining a BP value of the subject, based on each of successive pulses contained in the pressure signal, according to a relationship between blood pressure (BP) and pressure signal (PS). In the BP monitor, the BP-PS relationship is determined, and updated, based on the standard BP values periodically measured using the cuff. The cuff is necessary to measure the standard BP values used for calibrating the pressure signal of the surface-pressure sensor, because the accuracy of the BP measurement using the surface-pressure sensor is not satisfactorily high."}
{"text": "This invention relates generally to valves for molded inflatable articles, particularly to valves which accept an inflating needle, and to the method of forming the valve during the molding process.\nThere are several types of valves known for inflating articles such as balls, toys, tubes and similar articles. One type, represented by U.S. Pat. Nos. 3,107,683 and 4,341,382 utilizes metal valve parts inserted after the inflated article is formed. Another type, represented by U.S. Pat. Nos. 3,100,641, 3,204,959 and 3,220,729, utilizes metal or plastic parts which are placed in the mold prior to forming the inflated article. This valve is in the form of a tube or sleeve including a side aperture which pushes out or displaces the stem when pumping begins. Another type, more pertinent to this invention, is represented by U.S. Pat. Nos. 2,295,804, 2,302,985, 2,318,115, 2,387,455, 2,600,862, 2,760,775, 2,934,344, 3,174,501 and 4,660,831. These patents relate to inflatable articles in which the valves are self-sealing after withdrawal of an inflating needle. The valve chamber either contains a soft, sticky, puncture-sealing composition or a plug with a passage in the form of a slit which is effectively sealed by air pressure on the outside of the walls of the chamber which holds it. Inflation, thus, closes the plug opening. In each instance the plug is forced into the chamber by compressing it, so that it can be held in the chamber by some overlapping element such as a shoulder, a rib, a flange or a groove. The disadvantage of such an arrangement is that the pressure required to insert the plug, equals the pressure bringing about its removal. The use of a bonding agent for the plug is not totally satisfactory because it limits valve replacement.\nThe invention herein overcomes these disadvantages in a manner not disclosed in the known prior art."}
{"text": "Field of the Invention\nBelonging to a technical field of sunshade, the present invention relates to a sunshade, and more especially, relates to an opening and closing structure for an umbrella combining automatic and non-automatic modes applied to a sunshade.\nDescription of the Related Art\nAs modem daily appliances, umbrellas can roughly be divided into folded ones and straight ones. With the progress of the times, umbrellas with automatic opening and closing functions have been unveiled. With regard to the structure for an umbrella with automatic opening and closing functions, the design for the structure of a folded umbrella is easier to achieve due to the short stroke distance of its opening and closing movement, but for a straight umbrella with a long stroke distance of opening and closing movement, its structure is more difficult to achieve. An application which Publication Number is CN101999785A discloses an electric umbrella, including a movable rod and a tubular umbrella column. A lower end of the movable rod extends into the umbrella column and an upper end of the movable rod is fixedly connected with an upper umbrella plate, the umbrellacolumn is sleeved with a lower umbrella plate, and the upper umbrella plate and the lower umbrella plateare connected with an umbrella rib. An upper end of the umbrella column is fixedly provided with a pulley structures, the pulley structures are all connected with a rope, and the lower end of the movable rod is provided with a center hole, one end of the rope is connected and passes through the center hole, while the other end respectively fixed onto the lower umbrella plate. A transmission mechanism, a motor, and a control panel respectively connected with the motor are mounted in the umbrella column, wherein the transmission mechanism is fixed to the lower end of the movable rod and the motor's output shaft is connected with the transmission mechanism and capable of driving the movable rod to slide up and down through the transmission mechanism. However, an umbrella with this structure is difficult to manufacture due to the complexity of the structure and high costs. Nowadays, some electric sunshades are getting popular and are widely used in outdoor areas of venues such as beaches, villas, swimming pools and coffee shops, but the electric umbrellas are driven by a motor, so when there is no external power supply or no adequate battery charging, it is impossible to realize the electric opening and closing functions."}
{"text": "Traditional engine fuel control systems use a mechanical linkage to connect the accelerator pedal to the throttle valve. Engine idle speed is then controlled by a mechanical system that manipulates the pedal position according to engine load.\nSince the mid-1970's electronic throttle control or “drive-by-wire” systems have been developed. Electronic throttle control systems replace the mechanical linkage between the accelerator pedal and the throttle valve with an electronic linkage. These types of systems have become increasingly common on modern automobiles.\nGenerally, at least one sensor is typically placed at the base of the accelerator pedal and its position is communicated to the engine controller. At the engine, a throttle position sensor and an electronically controlled motor then regulate the throttle to maintain a precise engine speed through a feedback system between the throttle position sensor and the electronically controlled motor. An example of an electronic throttle control system can be found with reference to U.S. Pat. No. 6,289,874 to Keefover, the entire specification of which is incorporated herein by reference.\nIn conventional electronic throttle control systems, the various components of the throttle position sensor stator and connector assembly are mounted to the casting. The connector assembly is also connected to the motor. Thus, the throttle position sensor stator and the connector assembly move simultaneously during assembly and thermal expansion, thus possibly allowing one or both of them to become misaligned, which could potentially affect performance of the electronic throttle control system."}
{"text": "Webs of microfibers as taught in the prior art have significantly limited utility utility as thermal insulation. This is true despite the fact that previous investigators of microfibers have almost routinely included thermal insulation in their lists of potential uses for the fibers. See several patents dealing with blown microfibers, including Francis, U.S. Pat. No. 2,464,301; Ladisch, U.S. Pat. No. 2,571,457; Watson, U.S. Pat. No. 2,988,469; and Buntin et al, U.S. Pat. No. 3,849,241 (blown microfibers are very fine, discontinuous fibers prepared by extruding liquified fiber-forming material through orifices in a die into a high-velocity gaseous stream, where the extruded material is first attenuated by the gaseous stream and then solidifies as a mass of the fibers); or see Vinicki, U.S. Pat. No. 3,388,194, which describes microfibers formed by a centrifugal spinning operation. While these previous workers did not report thermal insulation values for a web of microfibers, or give any indication of having measured such values, they appear to have automatically considered that a new fibrous web should be useful as thermal insulation.\nThe limited value of existing microfiber webs as thermal insulation is true despite another fact, which is not set forth in any known prior literature, namely that microfiber webs provide unique insulating values. For example, a one-centimeter-thick web of polypropylene blown microfibers will give 1.8 clo of thermal resistance, in contrast to the about 0.9 exhibited by a 1-centimeter-thick web of commerical polyester staple fibers..sup.1 (Footnotes and methods of measurement are at the end of the specification).\nThe reason why existing microfiber webs have only limited value as thermal insulation is that, at least after these microfiber webs have undergone a normal compression history, they are heavier than alternative types of fibrous insulation. This heaviness is an inherent consequence of the very nature of microfibers; their very fine size and comformability causes the microfibers to come together as a dense, fine-pored web. As an illustration, a one-centimeter-thick web of blown microfibers is about five times as heavy as a one-centimeter-thick web of commercial polyester staple fibers. Even if a blown microfiber web only half as thick as a web of polyester staple fibers is used (so as to provide roughly equivalent thermal insulation resistance), the blown microfiber web will still be about two-and-one-half times as heavy as the polyester staple fiber web.\nWhen weight is only of secondary importance (as in glove or boot insulation, for example), thin, dense, blown microfiber webs can be quite useful. But when weight is a more primary consideration, as in such insulated articles as coats, snowmobile suits, sleeping bags, etc., existing microfiber webs will be passed by. Since the latter kinds of uses are major ones, the foreclosing of microfibers from such uses is a severe limitation on their utility in the insulation field."}
{"text": "In recent years, in a high-speed and high-density signal transmission between electronic devices or between wiring boards, the conventional transmission through electric wirings has began to reveal limitations in realization of high speed and high density, because mutual interference and signal attenuation constitute barriers. In order to overcome such limitations, a technology for connecting between the electronic devices or between the wiring boards by light, a so-called optical interconnection, is being examined. A polymer optical waveguide has drawn attention as an optical transmission path because of being easily processed with low costs, greater wiring-flexibility, and high density.\nAs the modes of a polymer optical wave waveguide, it is considered that a type which is formed on a glass epoxy resin substrate on an assumption of being applied onto a photoelectric consolidated substrate or a flexible type having no hard supporting substrate supposed to make a connection between boards are preferred.\nThe polymer optical waveguide is also required to have high heat resistance in addition to have high transparency (low optical transmission loss) in terms of usage environment of an applicable apparatus, part implementation thereof, or the like. Known materials for the optical wave guide include (meth)acrylate polymers (see, for example Patent Documents 1 and 2).\nHowever, even the (meth)acrylate polymers described in those patent documents have a high transparency of 0.3 dB/cm at a wavelength of 850 nm, evaluation of heat resistance performance, for example, specific test results such as optical transmission loss after a solder-reflow test are not described in detail and thus unclear.    Patent Document 1: JP Hei 06-258537 A    Patent Document 2: JP 2003-195079 A"}
{"text": "The Fischer-Tropsch (FT) synthesis has been used for long time for the production of synthetic hydrocarbons from synthesis gas (syngas), a mixture of gases comprising mostly hydrogen (H2) and carbon monoxide (CO).\nThe FT synthesis is a chemical reaction conducted over a metal oxide catalyst where the active metal comprises iron (Fe), cobalt (Co), ruthenium (Ru) and nickel (Ni). These catalysts can be produced by precipitation or can be supported in a metal oxide like alumina, titania, zirconia, magnesia and the like.\nThe FT synthesis primary reaction can be described as follows:2H2+CO→—[CH2]—+H2O\nIn this reaction —[CH2]— represents the primary building block of a paraffinic hydrocarbon, sometimes also referred as alkanes.\nThe process is carried out at elevated temperatures normally in the 200-400° C. range and high pressures of up to 40 bar-g. It can be conducted in many reactor designs like (i) tubular fixed beds, (ii) 3-phase slurry beds (a.k.a. bubble columns), (iii) high temperature circulating beds, and (iv) 2-phase fluidised beds. These had been extensively described in the literature in works like the Fischer-Tropsch Technology book edited by A P Steynberg and M E Dry (Elsevier, 2004).\nThis reaction proceeds by a mechanism known to those in the art as chain propagation. The reaction above described repeats many times resulting in the production of long chain species with up to 100 carbon atoms. The products from this synthesis are often described in the art as primary FT products, covering a very broad distillation range including FT hydrocarbon gasses (C2 to C4), FT liquid products (C5 to C21) and solid FT products (C22 and heavier) at ambient conditions. This blend has also been described as a synthetic crude or syncrude in many publications like “Processing of Fischer-Tropsch Syncrude and Benefits of Integrating its Products with Conventional Fuels” (NPRA Paper AM-00-51, 2000).\nThe primary products or syncrude from the FT synthesis can be separated into various liquid streams at processing conditions for example a C3 to C5 range, a naphtha C6 to C8 range and a C9 and heavier range.\nFuels refineries, irrespective of whether they refine crude oil, Fischer-Tropsch syncrude, coal liquids, oil shales or tar sands, generally produce a product slate that may include a kerosene cut.\nThe most common refining pathway for producing jet fuel (Jet A-1) from the kerosene range material is by hydroprocessing. In some instances it is desirable to produce more kerosene range material to ultimately boost jet fuel production and suggestions for doing this have been made, for example U.S. Pat. No. 4,409,092.\nMore recently, the United States Department of Defense expressed interest in a universal fuel for military use and preferably from a synthetic process to improve energy security (Forest, et al. 2005). This so-called Battlefield Use Fuel of the Future (BUFF), is very similar to jet fuel in specifications, but has a more stringent flash point specification (60° C. like JP-5). Ideally such a synthetic fuels refinery should produce only kerosene range material. Although a “jet fuel only” refinery may be conceptually devised, there are limits to the yield of kerosene range material that can be obtained in practise. Even conversion processes known in the art to be very kerosene selective, such as the conversion of propylene over solid phosphoric acid (Jones 1954), do not exclusively yield kerosene range material.\nStraight run Fischer-Tropsch products have some inherent drawbacks in meeting Jet A-1 and/or BUFF specifications, namely a high linearity that results in a high freezing point and low temperature viscosity and a tow aromatics content. Furthermore, the Anderson-Schultz-Flory distribution often used to describe the carbon number distribution of Fischer-Tropsch products, show that the volume of straight run syncrude material in the kerosene range is limited, irrespective of the Fischer-Tropsch process. Using the Anderson-Schultz-Flory description it may be shown that the straight run kerosene production from a Fischer-Tropsch process can be optimised by tailoring the Fischer-Tropsch catalyst to have a chain growth probability factor (α-value) preferably In the range 0.76-0.86. Most current commercial Fischer-Tropsch technologies operate outside of this range, with high temperature Fischer-Tropsch (HTFT) technology operating with an α-value of less than 0.7, while low temperature Fischer-Tropsch (LTFT) technology generally operates with an α-value higher than 0.9. However, even with a Fischer-Tropsch conversion process optimised for the production of kerosene range material, the straight run yield of kerosene is less than 30%. Nevertheless, the synthesis of jet fuel components from Fischer-Tropsch products and the blending thereof to produce semi-synthetic jet fuel and fully synthetic jet fuel are known in the art.\nIn this context, kerosene may be understood as a fraction having a carbon number range of C9 to C16, typically having a boiling point range of from 149° C. to 288° C., although variations to this definition may exist.\nPreparation of Jet A-1 from low temperature Fischer-Tropsch (LIFT) hydrocracker products by distillation results in a small C8-C13 fraction that meets most of the specifications, but contains less that the required 8% aromatics of Jet A-1\nThe initial boiling point of the material is limited by the flash point specification of jet fuel, while the final boiling point is restricted by the Jet A-1 specification (maximum 300° C.). However, in practise a lower final boiling point is often required to meet the freezing point specification (maximum −47° C.). It has also been pointed out that the low temperature viscosity of kerosene range material from Low Temperature Fischer-Tropsch (LTFT) syncrude may present a problem by being too viscous (Lamprecht 2006).\nThe LTFT hydrocracker design of Shell, as commercially used in their gas-to-liquids facility in Bintulu, Malaysia, has a maximum kerosene mode of operation which is reportedly capable of refining 50% of the LTFT syncrude to kerosene range material while the naphtha and distillate products from such an operation are only blending components and not transportation fuels meeting motor-gasoline and diesel fuel specifications.\nPreparation of Jet A-1 from high temperature Fischer-Tropsch (HTFT) hydrogenated straight run kerosene and iso-paraffinic kerosene from short chain olefin oligomerisation over Solid Phosphoric Acid (SPA), meets all the specifications, including density, however, the volumes that can be produced are unfortunately limited.\nWhat is needed is a synthetic fuels facility that can maximise kerosene yield in such a way that the kerosene meets Jet A-1/JP-8 and/or JP-5/BUFF specifications.\nA need has been identified for a kerosene dominant refinery in which the naphtha and distillate that are co-produced are easily refinable to meet final fuel specifications, rather than having to be sold as naphtha or distillate blend stock. This is especially pertinent in instances where the refinery is located far from markets for such intermediate products, or in instances where the refinery is used as strategic asset for the production of fuel for military use."}
{"text": "Mobile computing devices have been developed to increase the functionality that is made available to users in a mobile setting. For example, a user may interact with a mobile phone, tablet computer, or other mobile computing device to check email, surf the web, compose texts, interact with applications, and so on.\nBecause mobile computing devices are configured to be mobile, however, the devices may be exposed to a wide variety of environments having varying degrees of safety for the computing device. Accordingly, devices were developed to help protect the mobile computing devices from their environment. However, conventional techniques to install and remove the devices from the computing device alternated between being difficult to remove but providing good protection or being relatively easy to remove but providing limited protection."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of fabricating a semiconductor package, and more particularly, to a chip scale package (CSP) fabrication method including an electric die sort (EDS) and a burn-in test.\n2. Description of the Related Art\nSemiconductor packages have been continuously improved to meet new demands for high-performance, small-size and high-speed electronic appliances. Accordingly, semiconductor packages have been evolved from old-fashioned dual in-line packages (DIPs) to state-of-the-art semiconductor packages such as CSPs. As a result, electronic appliances such as camcorders or portable telephones can be miniaturized using the CSPs as their essential components.\nThe CSPs are known to be ideal for a package having 100 pins or more with high speed operation, or for a package having a large size chip. Although the definition of the CSPs has not been fixed in semiconductor industry, the packages that are smaller than 120% of the size of a semiconductor chip is typically referred to as a CSP. Even semiconductor packages larger than 120% of the size of a chip, such as ball grid array (BGA) packages, land grid array (LGA) packages, and small outline non-leaded (SON) packages, are considered as a CSP. In particular, the BGA packages having solder balls instead of leads mounted on the lower surface thereof, the LGA packages having a land array mounted on the lower surface thereof, and the SON packages having two land arrays instead of leads mounted on the lower surface thereof are examples of the CSPs.\nFIG. 1 is a flowchart illustrating a typical process for fabricating CSPs. Referring to FIG. 1, a plurality of semiconductor chips are formed on a wafer by integrated circuit fabricating processes. Then, the plurality of chips on the wafer are tested to identify good dies from bad dies depending on electrical characteristics or functionality of each chip by performing an electric die sorting (EDS) process. The EDS process utilizes a probe card having probe needles mounted thereon for electrically connecting a tester to an individual chip for testing, in step 51. Following the EDS process, preferably, only good dies are packaged into a CSP strip, in step 53. Next, the CSP strip is singulated (separated) into individual CSPs, in step 55.\nThe individual CSPs are subjected to a burn-in test for inducing the failure of marginal CSPs with defects, in step 57. This burn-in test accelerates failure mechanisms such that devices which have the potential to fail later but which failure would not otherwise be apparent at nominal test conditions can be eliminated by applying a serious stress, for example, heat, voltage or high frequency such as a clock, to a semiconductor package. Thus, defective CSPs can be screened out early before a final test, and only non-defective CSPs are subjected to the final test, in step 59.\nNext, the final test is preformed, in step 59, on CSPs that are determined to be good (non-defective) in the burn-in test. Then, CSPs determined to be good (fully-funictional) by the final test are marked their corresponding product titles, in step 61. Also, in step 61, the CSPs are finally inspected for visual defects. Thereafter, good CSPs are surface-mounted on a module board, in step 63. The CSP module is subjected to a module test, in step 65, and finally packed to be sent to end users, in step 67. Here, reference numeral 69 denotes a section where process steps are being carried out with singulated CSPs, before module assembly. Reference numeral 71 denotes a section from module assembly to final packing, where process steps are being carried out with a CSP module board.\nFIG. 2 is a schematic view illustrating a final test using a plastic tray 93 that is generally used to transport CSPs. Referring to FIG. 2, singulated CSPs are placed on the plastic tray 93 and then are transported between the processes using the tray 93. Thus, in the final test, singulated CSPs are contained on the test tray 93. The test tray 93 has a plastic main body 89 whose surface is treated as anti-static, and grooved pockets 91 on which singulated individual CSPs 81 are placed. Thereafter, the singulated CSPs placed on the test tray 93 are loaded onto a handler connected to a tester. In the handler, solder balls acting as external connection terminals for CSPs 81 are connected to socket contact terminals 87 (e.g., POGO pins) in a socket 85 on an interface board 83, and the CSP 81 is subjected to the final test. Although not shown, in burn-in test equipment, CSPs are inserted into sockets on a burn-in board, and subjected to a reliability test to sort out defective CSPs early.\nAccording to this conventional CSP fabrication method, singulated individual CSPs are loaded on a plastic tray and unloaded from a burn-in board numerous times during a burn-in-test. Also, in a final test, loading and unloading to and from the plastic tray are repeated.\nHence, the final test and the burn-in test with respect to singulated individual CSPs have the following problems.\nFirst, as external connection terminals, CSPs generally use solder balls, rather than outer leads. Such solder balls could be easily damaged by outside impact or physical stress. Further, encapsulating materials for the CSPs are relatively weaker than that of conventional encapsulating materials for ceramic or plastic packages. Thus, generally the CSPs are vulnerable against physical impact during the repeated loading and unloading in the burn-in test and the final test. Accordingly, various visual defects can occur easily. The visual defects include package breaking, cracking, damage to solder balls, and consequential contact failure. These visual defects significantly reduce the yield of the CSP fabrication process.\nSecond, the CSPs being a surface-mounting semiconductor package have more contact terminals than the other conventional semiconductor packages. Thus contact failure in the CSPs is more and more frequent due to an increase in the number of contact terminals during the burn-in test and the final test.\nThird, a complicated mechanism is needed to connect the socket to solder balls of the CSPs. This inevitably increases the price of socket, in tun, increasing overall manufacturing costs.\nFourth, the burn-in testing process and the final test require a great amount of time and man-power for repeated loading and unloading. Accordingly, an excessive amount of time and man-power is required in the conventional CSP fabrication process, thus lowering productivity.\nAccordingly, what is needed is a better way to fabricate the CSPs minimizing the visual defects, contact failure and the other problems discussed above, thereby improving productivity and lowering manufacturing costs."}
{"text": "Memory technology implies the need for accessing data within memory devices, such as DRAM or other memory devices. Therefore, different addressing strategies have already been proposed.\nA possible feature of DRAMs is address multiplexing. This technique enables splitting the address in half and feeding each half in turn to the chip on the address bus pins.\nThe chip has a large array of memory capacitors that are arranged in rows and columns. To read one location in the array, the control circuit first calculates its row number, which it places on the DRAM's address pins. It then toggles the row address select (RAS) pin, causing the DRAM to read the row address. Internally, the DRAM connects the selected row to a bank of amplifiers called sense amplifiers, which read the contents of all the capacitors in the row. The control circuit then places the column number of the desired location on the same address pins, and toggles the column address select (CAS) pin, causing the DRAM to read the column address. The DRAM uses this to select the output of the sense amplifier corresponding to the selected column. After a delay called the CAS access time, this output is presented to the outside world on the DRAM's data I/O pin.\nTo write data to the DRAM, the control logic uses the same two-step addressing method, but instead of reading the data from the chip at the end of the operation, it provides data to the chip at the start of the operation.\nAfter a read or write operation, the control circuit returns the RAS and CAS pins to their original states to ready the DRAM for its next operation. The DRAM requires a certain interval called the precharge interval between operations.\nOnce the control circuit has selected a particular row, it can select several columns in succession by placing different column addresses on the address pins, toggling CAS each time, while the DRAM keeps the same row activated. This is quicker than accessing each location using the full row-column procedure. This method is useful for retrieving microprocessor instructions, which tend to be stored at successive addresses in memory.\nIn addition, the provision of commands is typically provided through command strobes on the command bus. The number of different commands depends on the number of pins on the command bus with 2N commands being a possibility with N being the number of pins at the command bus. With the increasing demand for different commands, the command bus needed to be expanded. However, as die size is a crucial factor in application specific integrated circuit (ASIC) design, the number of pins on the command bus needs to be decreased.\nFurther, the size of the buses is also relevant for the overall size of the connection interface between the memory device and the central processing unit (CPU). The higher data rates that were required, the higher the number of connection pins on the data bus were selected. This increased the size of the interface. In addition, the number of pins on the address bus and the overall number of pins of the interface determined the type of memory to be used on the interface, besides protocol issues. However, the demand for flexibility of usage of different kinds of memory devices was not accounted for. There is a need for a flexible interface, which enables the use of different kinds of memory devices with different kinds of capabilities in terms of data throughput on the data bus.\nWith the rising need for flexible use of standard components, there is a need for providing a memory interface enabling to use both volatile memory and non-volatile memory on one same interface. However, as non-volatile memory and volatile memory have different prerequisites for the interfaces, there needs to be a possibility to adapt the interface for use with both types of memory."}
{"text": "The present invention relates to a newspaper conveying mechanism of a newspaper vending machine, which includes friction wheels and driving claw means to convey the newspaper at two stages. It is avoided that many newspapers are delivered at the same time. Also, the newspaper is prevented from being torn apart during conveying and is released as an integer.\nAs shown in FIG. 10, a conventional newspaper vending machine includes a newspaper conveying mechanism 1a, a newspaper releasing means 2a and an automatic changing mechanism 3a. The newspaper conveying mechanism 1a utilizes the weight of a main body to depress a pile of newspapers and employs rotary sharp claws 6a of a pushing means 5a driven by a motor 4a to forcedly directly contact with an upmost newspaper and push the same outside the machine.\nAccording to such arrangements, the sharp claws 6a directly contact with the surface of the newspaper and are slightly planted thereinto at several points for conveying the newspaper. Due to the pressure of the main body and the stress of the sharp claws 6a during rotation, the sharp claws 6a tend to damage the surface of the newspaper and imprint a track thereon. Moreover, in case that the motor 4a rotates at irregular speed or inappropriate speed, the newspaper may be torn apart. Therefore, it is necessary to provide an improved newspaper conveying mechanism to eliminate the above shortcomings."}
{"text": "Recent domestic events have heightened the need for an effective manner in which to assure the uncontaminated delivery of contained products to a consumer, particularly medicinal products taken internally. Specifically needed is the provision of a container for such products which bears assuring indication to the consumer that the contents have not been tampered with from their point of manufacture to the point of consumer sale.\nIn one prior art approach toward meeting this need, use is made of so-called \"telltale\" indication, i.e., a readily discernible characteristic indicative of tampering, such as a visible sign that some person has previously attempted to gain access to the container contents. Broadly speaking, these efforts may be generalized as placing a tamper-indicating member, e.g., an ambient-sensitive element, in the path of access to a container to indicate tampering by discernible change, e.g., change of color of the member. A quite early example of this practice is seen in U.S. Pat. No. 1,095,313 wherein a light-sensitive label is applied to a bottle and the releasably capped bottle with such label is wrapped in a light-impermeable paper. When the wrapper is removed in ambient light, the label changes color and indication is thus provided to a subsequent purchaser that the wrapper has previously been removed. In a practice within the last decade, seen in U.S. Pat. No. 3,899,295, this technique is modernized by including the telltale substance as an interiorly disposed protected component of the wrapper. In the '295 patent, a heat-shrinkable member straddles both the cap and container vessel after capping and has a pH-sensitive integrity indicia imprinted on the interior of the member, the indicia being packaged with a basic gaseous material which maintains the indicia of a given first color. When the heat-shrinkable member is first removed from the cap and vessel, ambient pH causes the indicia to change color.\nAnother telltale approach is seen in situations in which containers are not releasably capped, i.e., the telltale is a component of a strippable closure member. Examples of this effort are seen in U.S. Pat. No. 3,826,221 and 3,923,198. In the '198 patent, a multilayer member serves to close the access avenue to a container and includes a layer which becomes opaque when subjected to stress. A color backing or printed legend normally visible through the stress-sensitive layer is not seen on tampering, thus providing a color change which is discernible to the user to indicate that tampering has occurred. In the '221 patent, an outer seal is adhesively secured to a container as a closure member and includes an ink which smudges if the closure member is tampered with.\nIn applicants' view, the latter approach is more desirable in one aspect than the former, since the latter provides indication of tampering directly at the access port rather than at a preceding wrapper removal stage. Thus, the heat-shrinkable member discussed above is a stage removed from the removal of the cap of the container and may not be present at the cap removal. However, such advantage in the latter techniques is obtained at the expense of exposing the telltale to inadvertent activating stress in the course of handling and shipping. It is applicant's further view that tamper indication should be effectively provided without need for such ambient-sensitive telltales or that more effective such ambient-sensitive telltale containers should be afforded to manufacturers."}
{"text": "Wireless communication systems are constantly evolving. System designers are continually developing greater numbers of features for both service providers as well as for the end users. In the area of wireless phone systems, cellular based phone systems have advanced tremendously in recent years. Wireless phone systems are available based on a variety of modulation techniques and are capable of using a number of allocated frequency bands. Available modulation schemes include analog FM and digital modulation schemes using Time Division Multiple Access (TDMA) or Code Division Multiple Access (CDMA). Each scheme has inherent advantages and disadvantages relating to system architecture, frequency reuse, and communications quality. However, the features the manufacturer offers to the service provider and which the service provider offers to the consumer are similar between the different wireless systems.\nRegardless of the modulation scheme in use, the wireless phone available to the end user has a number of important features. Nearly all wireless phones incorporate at least a keyboard for entering numbers and text, and a display that allows the user to display text, dialed numbers, pictures and incoming caller numbers. Additionally, wireless phones may incorporate electronic phonebooks, speed dialing, single button voicemail access, and messaging capabilities, such as e-mail.\nSubscribers' reliance upon their mobile terminals (such as, cell phones) and on wireless service in general continues to increase. The demand for and implementation of Emergency-911 and geo-location services for mobile terminals are two typical examples of this increasing demand. However, subscribers quickly learn that, once away from areas equipped with cells and base stations, their mobile terminals are no longer able to function. This lack of connectivity is a significant drawback of the prior art."}
{"text": "1. Field of the Invention\nThe invention relates to the field of power electronics,wherein mains system disturbances and mains voltage reductions occur in a voltage supply system.\n2. Discussion of Background\nCircuit arrangements of this type, also known as Dynamic Voltage Restorers (DVR), supply a separate voltage, in series with the disturbed voltage supply mains system, in order to compensate for mains voltage reductions and mains system disturbances. In consequence, only insignificant changes and disturbances can be recorded on the load. To this end, the feeder converter must react very quickly, which means that the power semiconductor switches in the converter have to have a high clock frequency. This also results in the requirement that the filter connected on the mains system side of the feeder transformer must be kept as small as possible. In conventional systems for low and medium power levels, it is possible to use IGBTs and a comparatively low intermediate-circuit voltage. During operation at the required clock frequencies in such configurations, the switched-on losses of the elements are typically the governing factor, to such an extent that no further measures can be taken to reduce the overall losses significantly. At higher power levels, it is on the one hand increasingly more difficult to use IGBTs, and on the other hand the intermediate-circuit voltage must be increased in order that the currents that occur can still be coped with well. In addition, conventional GTOs have the disadvantage that they cannot be switched fast enough. In general, the switching losses in silicon semiconductor switches tend to increase at higher voltage loads. This applies to both IGBTs and GTOs. Modern, hard-driven GCTs (Gate Commutated Thyristors) can, furthermore, be produced with lower switched-on losses than IGBTs. The ratio of the switched-on losses to the switching losses is thus shifted to a major extent toward the switching losses for high-power applications, particularly with respect to the use of GCTs."}
{"text": "The present invention relates to hair thickening and conditioning compositions, and in particular to aqueous gelatin solutions, the application of which to the hair results in a significant increase in hair shaft thickness. The present invention also relates to methods for preparing the gelatin compositions of the present invention, as well as to methods for increasing hair shaft thickness by topical application to the hair of the gelatin compositions of the present invention.\nU.S. Pat. No. 4,749,684 discloses that the ingestion of gelatin promotes an increase in the linear growth rate of hair. Scala et al., Nutr. Rep. Int., 13(6), 579-92 (1976) discloses that dietary supplementation with gelatin produces an increase in hair thickness. The ability to increase hair thickness by the topical application of gelatin is unreported.\nHair coated with gelatin is difficult to comb in the absence of a conditioning agent, but stable solutions of gelatin with a conditioning agent have been heretofore unknown. Lindo et al., \"Hair Processing and Conditioning,\" Hum. Hair Sym. [Pap.], 1st, (Brown, Editor, MEDCOM Press, New York 1974 (Meeting Date 1973)) pages 135-44 reports the capability of hair to absorb topically applied hydrolyzed gelatin, but with little substantivity, so that little gelatin remains after thorough washing with room temperature water.\nU.S. Pat. No. 4,208,402 discloses a topical solution for improving the appearance of oily hair or skin, including the scalp, containing from about 0.5 to 1.5% by weight of gelatin and from 0.15 to 0.45% by weight of partially hydrolyzed polyvinyl acetate in aqueous or hydroalcoholic solution. The level of film-forming polymer disclosed is insufficient to provide a readily combable product. By the manufacturing process disclosed, aqueous solutions cannot be formed with higher concentrations of the film-forming polymer without the gelatin precipitating from the solution.\nTopical application of gelatin to hair would be possible if stable aqueous solutions containing higher concentrations of film-forming polymers could be prepared. The inability to comb hair to which gelatin solutions have been topically applied has prevented the recognition that topical gelatin application increases hair shaft thickness.\nThere remains a need for stable aqueous gelatin solutions containing higher levels of film-forming polymer to provide gelatin solutions for topical hair application that do not resist combing."}
{"text": "1. Field of the Invention\nThe present invention generally relates to high voltage on-chip circuitry for electrically erasable programmable read-only memory EEPROMs or flash memory. More particularly a low power, high-voltage ramp-up rate control circuit is described which utilizes a current reference, a current controlled oscillator, a charge pump, a feedback transistor, a reference transistor, and a ramp-up transistor.\n2. Description of the Prior Art\nFIG. 1 contains a circuit schematic which shows the prior art implementation of a ramp-up rate control circuit. A functional block labeled charge pump+oscillator+current reference 110 is shown. The output of this block 110 is a charge pump current, Icp 170. In addition, the voltage on the output node 195 of this block 110 is HV. (High Voltage). A reference current, Iref, 120 is developed in a current source. This current flows into the drain of an N-channel metal oxide semiconductor field effect transistor, NMOS FET, M1130. The gate 175 of FET M1130 is connected to its drain. The source of FET M1 is connected to ground 185.\nThe ramp current Ir 190 flows into the drain 165 of NMOS FET M2140. The gate of FET M2140 is connected in common to the gate of FET M1175. The source of FET M2 is connected to ground 185. The drain 165 of FET M2 is connected to the gate of NMOS FET M3150. The drain 180 of FET M3150 receives the leaking current Ileak. The source of FET M3 is tied to ground 185.\nIn this prior art circuit in FIG. 1, the leak transistor M3150 will tend to discharge the charge pump output node HV 195. The leakage current Ileak 180 depends on the gate voltage V1165 of FET M3150. The ramp-up rate (dV/dt) is equal to Ir/C where C is the capacitance of the capacitor 160 and V is the voltage across capacitor C. In FIG. 1, the leakage path will consume the extra current Ileak. If this extra leakage path 180 can be eliminated, then the power consumption can also be reduced.\nU.S. Pat. No. 4,488,060 (Simko) “High Voltage Ramp Rate Control Systems” describes high voltage ramp rate control systems. These systems are useful for on-chip EEPROM high voltage power supplies.\nU.S. Pat. No. 5,872,733 (Buti et al.) “Ramp-Up Rate Control Circuit for Flash Memory Charge Pump” describes a ramp-up rate control circuit. This circuit controls the ramp-up rate of a charge pump having an output which provides an output voltage and an output current.\nU.S. Pat. No. 5,408,133 (Honnigford et al.) “Ramp Control Circuit” discloses an apparatus for providing an EEPROM programming signal.\nU.S. Pat. No. 5,365,121 (Morton et al.) “Charge Pump with Controlled Ramp Rate” discloses a charge pump with a controlled ramp rate. It is made up of a charge pump, a differentiator circuit and a trigger circuit."}
{"text": "Touch sensing and multitouch sensing are key technologies in the implementation of sophisticated modern human-machine interfaces. Touch sensing can involve sensing the proximity, contact, and/or position of an object such as, for example, a finger, stylus or other object. Multitouch sensing can involve similar sensing with respect to multiple simultaneous input objects. As such, multitouch gestures are now being implemented in almost every electronic device that has a touch interface. Typical touch sensing and multitouch sensing systems are based on measures of absolute or mutual capacitance. One example of multitouch sensing is multitouch imaging, which can involve capturing transcapacitive images relative to an input interface/sensor."}
{"text": "1. Field of the Invention\nThe present invention relates generally to systems and methods of providing synchronized information such as narratives, translations or other show related messages to individual patrons of rides, shows and exhibits.\n2. General Background and State of the Art\nThere are several circumstances that may interfere with a person\"\"s ability to hear, and thereby interfere with a presentation of information to the listener. For example, persons with hearing loss may miss narratives, sound effects, music and other sound material related to the presentation, and messages in live performances, films, television and special events. Persons who do not speak the language or languages used in the presentation may miss narratives and other related messages that are presented. The language barrier prevents many people from different cultures and languages from understanding, participating or interacting with the information being presented. Background environmental noise may also affect a person\"\"s ability to hear, and thereby diminish the effect of the presentation. Persons with severely impaired vision also miss visual portions of these presentations.\nCaptioning systems have been used in venues including museums, theaters and other auditoriums to provide foreign language translation or captioning for the hearing impaired. The most common of these systems includes xe2x80x9copen captioningxe2x80x9d whereby text is displayed on a projection surface or large adjacent display area where the entire audience can see the captioning. Another common method of captioning is xe2x80x9creflective captioningxe2x80x9d which uses a transparent but reflective panel to display the text from a rear projection while allowing the viewer to see the display or performance through the panel. Yet another method utilizes individual hard-wired displays located in the back of a seat in front of the viewer.\nMultiple channel assistive listening systems have also been used in venues including theaters, churches and auditoriums to provide amplified narration and foreign language translation. These systems have the content broadcast by infrared, radio frequency or inductive loops to the devices carried by the patrons. In these cases, the content is continuously streamed with the content source being located in the facility; the receiver simply played back the audio as it was received. In those instances where more than one language is provided, the multiple transmission frequencies sometimes cause interference with each other (crosstalk). The available number of suitable frequencies also limits the number of channels. Radio frequency and inductive loop technologies can also cause interference with other electronic equipment within the venue. Personal audio devices have also been used in venues including museums to provide narration in multiple languages. These are typically manually activated systems where the patron is required to activate the system, either by pressing a button or manually entering a numeric code.\nEach of these systems has its own limitations. For example, open captioning can be obtrusive and distracting to other patrons. These systems are generally not convenient to the user and are limited to one or a few languages.\nMoving ride vehicles in a venue such as an amusement park present an additional set of problems in providing captioning due to such factors as variations in vehicle speed, and lack of communications wiring to moving vehicles.\nAdditionally, presentations often could be made more effective if they included a method to provide interactivity between the audience and the presenter(s).\nIt is therefore an object of the present invention to provide a convenient system and method of presenting audio and/or visual information for the purposes of captioning, language translation, assistive listening, and descriptive audio to individual patrons of rides, meeting areas, pathways, shows and exhibits without using intrusive open captioning or hard wired individual display systems. The present invention preferably makes use of infrared (IR) or radio frequency (RF) signals to wirelessly trigger portable devices.\nWireless transmitters are installed at various locations including waypoints, signs, display areas, entrances, exits, and theater seating areas. The transmitters are configured to send short messages, or codes, corresponding to waypoint location, type of sign, show time, time of day, etc. The transmitters may be freestanding transmitters that repeatedly transmit stored messages, or the transmitters may send data from other equipment such as computers, show controllers, devices with processors or controllers, other transmitters, etc. The transmitters may cover a large area such as a theater, or be precisely directed to cover a small area.\nGuests carry portable devices loaded with content information such as location descriptions, sign contents, show content, alternate language content, event times, etc. Content may be audio, graphical, or a combination of both.\nWhen the device is within range of a transmitter sending codes, the device should be capable of receiving codes. Upon receiving a code, the device searches its memory for appropriate content, tests any logic constraints, and presents content to the user. Logic can be based on user preferences, time of day, time of event, locations of transmitters visited, time of visits, sequence of transmitters visited, intended device type, etc.\nAll portable devices have the common elements of receiver, processor, memory, and power source. Other elements of the devices vary depending on intended applications. A portable captioning unit may have a display. A display, as used herein, may be defined as a visual display for text or graphics, or for audio information. For example, a portable audio device has an audio amplifier with speaker or earpiece. A portable captioning unit may have a visual display. Other devices may have both audio capability as well as a visual display, and further may have user input capabilities. Portable devices may also have transmitters for communication with other portable devices, with system docking stations, with receivers in shows or attractions, etc. Portable devices may have IR receivers, RF receivers, or both. Portable devices similarly may have IR transmitters, RF transmitters, or both.\nContent information may be loaded into the portable device memory at any time prior to its use. A bulk storage device may be attached to the portable device, the portable device memory contents may be received via a docking station data interface, or content may be streamed into portable device memory via its wireless receiver. Examples of bulk storage are flash memory cards, micro-drives, hard drives, CD-ROM disks, etc.\nContent may be loaded into a portable device by connecting a cable, or plugging the device into a cradle or docking station having a connector. Examples of docking station data interfaces are Ethernet-based serial, RS-232 based serial, IRDA infrared serial, etc. It is possible to provide docking stations at numerous locations available to the portable device user. It is also possible to provide large, multi-unit docking stations that load many portable devices simultaneously. A docking station may connect to a freestanding storage device, or may be linked to one or many remote content sources via a communications network. The network could include a local area network (LAN), Internet connection, satellite link, etc. The content could also be streamed to a device using a wireless network. For example, streamed content could be loaded at a theater to send blocks of caption text to portable units shortly before sending text synchronizing commands; the content is buffered in portable device memory shortly before its presentation.\nOne such system could use high-speed RF transceivers to transfer content to portable devices, while using an IR system to precisely aim location synchronizing codes to the portable device. It is also possible for a portable device to automatically request content download; as a user travels to a new area with a portable device and receives location codes for which it has no content, the device may transmit a request for new content. Such a content delivery system may connect to a freestanding storage device, or may be linked to one or many remote content sources via a communications network. The network could include a local area network (LAN), Internet connection, satellite link, etc.\nIn one embodiment of the present invention, text captions are synchronized with a theater presentation. A portable device with a text display receives time codes from a theater transmitter. Caption text for the entire theater presentation is stored in the device. The device receives the transmitter signals, extracts the current show time, and displays the appropriate text in time with the show.\nIn another embodiment of the present invention, a personal audio device receives location codes from transmitters mounted along a travel route. Transmitters are located at significant locations along the travel route; transmission coverage at each location is restricted. As the user nears a location, the device receives the transmitter signals, searches its memory for matching audio content, and plays the audio.\nIn yet a further embodiment of the present invention, a device plays speech translated into the user\"\"s language. The theater transmitter sends time codes during the theater presentation. The device memory contains show speech content recorded in the user\"\"s language. The receiver in the device extracts the current show time, and plays the matching speech. The received time codes will synchronize show speech even if the user has entered the theater mid-show.\nThe system and method of the present invention is designed for use in a wide variety of locations. For example, the present invention is suitable for a theme park that contains a large number of different attractions covering a large area and receives a large number of visitors, including tourists who may not speak the native language or visitors with specials needs due to audio or visual impairment. The invention could also be used in theaters, museums, shopping malls, exhibit halls, convention centers, meeting areas, walkways or roadways, and a wide variety of other locations."}
{"text": "Such a transformer and such a lamp base are disclosed, for example, in the laid-open specification WO 00/59269. This specification describes a toroidal core transformer, whose windings are surrounded by a housing, and which is arranged in a chamber of the lamp base of a high-pressure discharge lamp."}
{"text": "1. Field of the Invention\nThis invention relates to inductor type dynamoelectric motors and, more specifically, to a support and positioning structure for conductors of a stator assembly of such a motor.\n2. Background Art\nDynamoelectric machines are commonly used as motors for converting electrical to mechanical energy. Inductor type dynamoelectric motors may be used as adjustable speed drives for high speed operation. A dynamoelectric machine may also be used for electrical generation by driving its rotor in reverse thereby producing energy in the armature windings. Inductor type dynamoelectric machines are generally characterized by a stator which includes both AC armature and DC exitation coils surrounding a coil-less rotor. In this configuration, there are no rotating field or armature coils, slip rings, brushes and associated connections which are common to machines having rotating windings. Since inductor type dynamoelectric machines contain fewer rotating parts, have a more rugged rotor construction and are particularly suitable for high speed application. One version of such an inductor type dynamoelectric machine, employs a circumferentially distributed arrangement of C-shaped armature elements surrounding a generally cylindrical field coil which in turn encloses a transverse pole magnetic rotor. U.S. Pat. Nos. 437,501 and 2,519,697 and 3,912,958 describe earlier machines of this type of design.\nMore recent versions of such inductor type machines are disclosed in commonly owned U.S. Pat. Nos. 4,786,834 and 4,864,176. In these patents, a spool-like support structure supporting field windings and armature elements is disclosed. The spool-like structure is made of a nonmagnetic material and has a hollow, elongated central portion extending concentrically about a longitudinal axis. This central portion supports a field coil and defines an interior longitudinal passageway for accommodating the insertion of a coaxial rotor. End portions are located at each end of the central portion and extend radially outward therefrom. Each of the end portions is preferably provided with radially oriented grooves in its axially outermost surface. The grooves are configured to receive and orient legs of generally C-shaped armature coil elements arranged in a circumferentially distributed pattern about the periphery of the spool-like structure. The armature elements are thereby positioned in three orthogonal directions.\nOther features, advantages and benefits of these dynamoelectric machines, including the stator support structures thereof, are described in detail in U.S. Pat. Nos. 4,786,834 and in 4,864,176. Each of these patents is incorporated herein by reference and made a part of this disclosure.\nThe C-shaped armature elements contain armature windings thereon which carry an AC current and generate a magnetic flux which drives the rotor when the machine is used as a motor. In this situation, the current must be supplied to each of the circumferentially oriented C-shaped armature element's windings. Typically, low power, dynamoelectric machines require less current through the armature windings thereby requiring lighter gauge windings. Accordingly, the conductors transmitting current to the windings may also be of a relatively light gauge. Because lighter gauge conductors are used in these low power, dynamoelectric machines, there have been few problems connecting the armature windings to the conductors, supporting the conductors, and connecting the conductors from the windings to a power supply wire outside of the dynamoelectric machine. For example, a light gauge power supply wire is typically fed through the housing of a low power dynamoelectric machine. The wire is then connected to a printed circuit board which contains a network of conductors therein which distributes the current therethrough and into the armature windings which are typically soldered to portions of the network. In this situation, the conductors are of relatively light gauge and are supported within the circuit board.\nIn high power high speed rotation dynamoelectric machines, the power to supplied to the armature windings is greater than in low power machines. These machines require high frequency currents. However, at high frequency, the effective AC resistance of conductors increases with increasing frequency. To reduce the increase of conductor resistance at high frequency, the armature conductors may be multiple isolated ribbon-like strands. Therefore, heavier gauge conductors are needed to supply the necessary current. However, these higher gauge conductors are often too heavy and bulky to be supported by a printed circuit board or similar type support structure. Also, the conductors are not easily connected to the armature windings.\nIt is therefore desirable to provide a system for supporting conductors which supply current to the armature windings of a high power, high speed dynamoelectric machine. It is also desireable to provide a system for connecting the conductors to the armature windings and to provide a system which does not significantly increase the size of a dynamoelectric machine."}
{"text": "Conventional speech recognition systems attempt to convert speech from a user into text that represents the words that were actually spoken by the user. Such speech recognition systems typically suffer from a number of problems and are relatively inaccurate. That is, conventional speech recognition systems are often unable to determine what was spoken or incorrectly determine what was spoken. In either case, such problems limit the usefulness of these speech recognition systems."}
{"text": "Several service providers, such as those providing multicast services, e.g., Video over Internet Protocol (VoIP), IP television (IPTV), etc., deliver multicast content from one or more head-end nodes (e.g., Points of Presence, or POPs) to one or more tail end nodes (e.g., several tail-end nodes per head-end node) over a point-to-multipoint (P2MP) tunnel. In one scenario, the head-end node receives the multicast traffic from a multicast network and transmits the multicast traffic onto the P2MP tunnel. Also, the tail-end node receives the tunneled traffic, and forwards the received traffic to another multicast network located at the tail-end node, for example, to reach end-user devices (client televisions, etc.). Notably, the P2MP tunnel, e.g., established using Multiprotocol Label Switching (MPLS) Traffic Engineering (TE), offers a tunnel connection through the service provider's (and others') network, which, as those skilled in the art will understand, allows for various tunnel advantages to be achieved through the network, such as bandwidth reservation, etc.\nOne particular benefit of a tunnel is its ability to provide “Fast Reroute” (FRR) functionality to protect against intermediate node failure along a primary tunnel. That is, an FRR backup tunnel may be established to protect one or more nodes (or links) along the primary tunnel and, in the event the node (or link) fails, the point of local repair (PLR) quickly reroutes the primary tunnel traffic onto the backup tunnel to circumvent the failed element. However, because the FRR backup tunnel generally needs to intersect the primary tunnel (i.e., the backup tunnel generally starts and ends at the primary tunnel), failure protection is generally not available to a head-end node and a tail-end node of the tunnel.\nFor instance, for P2MP tunnels receiving multicast traffic, in the event that a head-end node of a primary P2MP tunnel fails, then a backup head-end node is configured to rerouted the multicast traffic onto a corresponding backup P2MP tunnel. This process may require cumbersome configuration (e.g., manual configuration) to determine the backup head-end node that may be suitable to assume the responsibility of the failed head-end node. Particularly, the backup P2MP tunnel has generally already been established (i.e., prior to failure) from the backup head-end node to each of one or more tail-end nodes originally receiving tunnel traffic from the failed head-end node. In addition, the primary and backup (and other possible backup) P2MP tunnels from the multicast network to the tail-end nodes inefficiently reserve resources (e.g., bandwidth) since the tunnels are computed separately, without consideration for shared resources (i.e., where the primary and backup tunnels utilize the same links/nodes)."}
{"text": "As is well known to those skilled in the art, alcohols such as ethanol may corrode metal surfaces with which they come into contact. This is particularly true of crude or commercially available ethanols which undesirably contain acidic components commonly acetic acid. In the case of fermentation alcohols, acetic acid may be present in amount of 0.003 w %-0.005 w % of the alcohol; and this may be responsible for the fact that the alcohol causes serious corrosion problems.\nIt is an object of this invention to provide a novel process for decreasing the corrosion of alcohol compositions. Other objects will be apparent to those skilled in the art."}
{"text": "The present invention generally pertains to systems and methods for accessing a hosted service over a network. More specifically, the present invention deals with methods for providing secure access to a hosted service via a client application.\nThe functionality of certain software applications can be extended through services offered through a network such as the Internet. Communication with the provider of services should be secure in order to protect the interests of both the host of the application and the service provider.\nRetail management systems are a specific area where securing communication with a remote service provider is challenging. In order to achieve some degree of automation, retail businesses often implement a specialized software application. Many of these applications are point-of-sale solutions that enable at least partial automation of any of a number of processes such as customer tracking and inventory management. One example of such a software application is Microsoft Retail Management System (MRMS) provided by Microsoft Corporation of Redmond, Wash. Other examples of such software applications include back office systems, store room and shipping applications, MRMS Headquarters and warehouse management software.\nIt is common for retail management software applications to be installed on multiple computers (e.g., connected by a Local Access Network) that operate in conjunction with a central database. In some instances, extended functionality is available to the retail application in the form of remote services delivered by a service provider through the Internet. Such extended functionality may include, by way of example, payment card processing, integration with e-commerce web hosting or merchandising services. These and other services may be provided for free or based on a payment scheme involving, for example, subscription or per access based charges such as billing per transaction and metered billing (e.g. based on disk usage, quality/speed/level of service).\nUser access is an important area of consideration for many of the described and other remote service systems. For example, distributing appropriate access rights to different users in some customized manner (e.g., different employees or employee roles are assigned different access rights) is often a desirable capability.\nSome hosted web services are only designed to support a single user login account per application account. This can be impractical in many environments, such as a retail sales environment wherein there is often a high turnover in staff and a need to provide access to multiple users (e.g. more than one person doing shipping of product sold on-line, different users on separate shifts, more than one person needed to update e-commerce website product listings). Furthermore, it is conceivable for a software application to provide its own user authentication system that eliminates the necessity of user authentication with a hosted service. It is desirable to provide “seamless” integration of an application and a hosted web service without requiring unnecessary log-in steps and password transactions. For example, it is undesirable to maintain and update separate employee user accounts for an on-line service."}
{"text": "In general, open-roof designs for individualized transportation have been utilized for many years as established by carriages and the first engine-driven vehicles. Open driving was and is part of the culture of automobile driving with continued interest for contemporary designs. In the case of convertibles, it is primarily folding roofs, which are constructed in such a manner as to be able to be retracted and stretched out again, with a structure of bars and an external skin which is correspondingly foldable. Such types of folding roofs have disadvantages, particularly for reason of restricted suitability during bad weather and winter and are, due to construction, adversely affected by unpleasant wind noises at high speeds.\nFor increased suitability during bad weather and winter, rigid convertible folding tops, so-called hard tops were designed to provide a detachable roof, with the same paintwork and finish as the body. However, early designs were extremely expensive and complex, lending to limited production offerings.\nBecause of tendencies within the culture of motor vehicles at the present time, convertibles are no longer purely used as leisure-time vehicles, such as a second or third household car. Instead, such vehicles serve as a primary vehicle, leading to significant design options for hard tops which can be retracted into the vehicle space behind the passenger compartment. The most common areas for storage include a space between the passenger compartment and the trunk or within the trunk itself, thereby restricting use of the trunk space while retracted. This design is commonly referred to as an internal storage design. Alternative designs have also included storage of the hard top above the vehicle trunk for purposes of simplifying the design and maintaining a pleasing appearance. Such designs are commonly referred to as an external storage design. However, such designs impede or prevent use of the vehicle trunk while the hard top is retracted and noticeably change the appearance of the vehicle profile. Therefore, a need exists for providing a retracting system that will allow external storage of the vehicle hard top that will not significantly affect the side profile of the vehicle and permit full use of the vehicle trunk space.\nThe present disclosure provides a solution to this and other problems known in the art, and offers other advantages over the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to the manufacture of pharmaceutical and pharmaceutical-like product. More particularly, the present invention relates to an apparatus for manufacturing pharmaceutical and pharmaceutical-like product.\n2. Description of Related Art\nContemporary quality control methods for pharmaceutical and pharmaceutical-like product involve the use of batch sampling techniques. The batch-sampling techniques test samples from batches of the product, such as through the use of wet chemistry, after the product has been made. Contemporary batch sampling techniques use frequent and sometimes random batch sampling for various characteristics of the final product, such as, for example, quality, concentration and homogeneity. However, these batch-sampling techniques suffer from drawbacks because of their inefficiency and inaccuracy.\nBatch-sampling assumes that all of the product attributes in a particular batch are normally distributed and have the same or very similar characteristics as the sampled product from the batch. Where the chosen samples do not meet the required tolerances, an entire batch can be discarded or re-processed for additional sampling and testing. If the chosen unacceptable samples do not have the same characteristics as other acceptable product in the batch, then acceptable product may be discarded along with the rejected samples or at least need to undergo more costly testing. Batch-sampling can be particularly inaccurate where the error or flaw in the process is random, non-repeating or of a non-linear nature. Such flaws or errors in the manufacturing process may provide for only a fraction of the product of the batch being unacceptable but result in an entire batch being discarded or re-tested, as a result of the use of batch sampling.\nAnother significant drawback of batch-sampling techniques is where the chosen samples meet the required tolerances, but where a fraction of the batch is actually unacceptable and not representative of the tested sample. In such a situation, unacceptable product may be provided to the consumer because of the inherent flaw in the quality control method.\nAn additional drawback in batch-sampling techniques is that the testing is done at the end of the process and provides little, if any, information for corrective action to be taken with regard to the manufacturing process and its various steps. The batch-sampling technique can provide overall information for sampled product, but does not indicate at which point or which particular step in the process that a flaw is occurring, such as, for example, inadequate dosing or detrimental heating.\nAnother drawback of batch-sampling technique is that it is done off-line of the manufacturing process, which adds time to the overall manufacturing process, and can also be labor intensive. The cost in time and labor is increased where more stringent standards are applied to a particular product so the batch-sampling technique utilizes a higher portion of samples for testing.\nAccordingly, there is a need for an apparatus and process for manufacturing pharmaceutical and pharmaceutical-like product that reduce or eliminate these manufacturing and quality control drawbacks of the contemporary devices and techniques."}
{"text": "The present invention relates to a conditioner, in particular for the production of conditioned foodstuffs or animal foods.\nConditioners are used individually as processors or in conjunction with pressing or expander plants in the preparation of animal foodstuffs or foodstuffs. These are generally mixers, in which a product for processing is processed over a certain period thermally or hydrothermally. This processing serves to render the products hygienic, and to improve the nutritional value and/or the product structure.\nSuch mixers are known for example from EP-B-231764 and EP-B-610789. In this example a dosing device is connected upstream of such a mixer and a compression mould or an expander is supplied from the mixer. The mixer is provided with a heating mantle and contains a rotating shaft, which is equipped with mixing and conveying elements.\nSuch units are either tailored to use in pressing or expander plants, or can be used stand-alone only, since the expense of adapting them is high.\nAccording to EP-B-610789, for warming the product-guiding parts and units, an air cycle is provided with a ventilator, a connecting line, an air heater, a lower air connection arranged on the side of the product output channel with an air inlet channel, an air input channel with an upper air connection arranged on the side of the product supply channel, whereby the air is guided from the product output channel via the compression mould, the product supply channel, the feeder and mixer to the products input channel.\nSolutions to energy recovery are also known in processors, in which for example the waste heat of an expander plant is utilized."}
{"text": "The present invention relates to poly[alkylene-4,4'-(ethylenedioxy)bis benzoate] copolymers as well as surgical devices formed therefrom. More particularly, this invention relates to flexible monofilament surgical sutures having unique handling and knot-tying characteristics.\nMany natural and synthetic materials are presently used as surgical sutures. These materials may be used as single filament strands, i.e. monofilament sutures, or as multifilament strands in a braided, twisted or other multifilament construction. Natural materials such as silk, cotton, linen, and the like, do not lend themselves to the fabrication of monofilament sutures and are accordingly used mostly in one of the multifilament constructions.\nCertain synthetic materials which are extruded in continuous lengths can be used in monofilament form. Common synthetic monofilament sutures include polypropylene, polyethylene and nylon 6. Such monofilament sutures are preferred by surgeons for many surgical applications due to their inherent smoothness and noncapillarity to body fluids.\nAvailable synthetic monofilament sutures all suffer to a greater or lesser degree from one particular disadvantage, that is relative stiffness. Besides making the material more difficult to handle and use, suture stiffness or low compliance can adversely affect knot-tying ability and knot security. It is because of the inherent stiffness of available monofilament sutures that many suture materials are braided or have other multifilament constructions with better handling, flexibility and conformity.\nMost monofilament sutures of the prior art are also characterized by a high degree of stiffness. This makes knot-tying difficult and reduces knot security. In addition, the low compliance and limited ductility prevent the suture from \"giving\" as a newly sutured wound swells, with the result that the suture may place the wound tissue under greater tension than is desirable, and may even cause some tearing, cutting or necrosis of the tissue.\nThe problems associated with the use of low compliance sutures in certain applications were recognized in U.S. Pat. No. 3,454,011, where it was proposed to fabricate a surgical suture composed of Spandex polyurethane. Such sutures, however, were too elastic and did not find general acceptance in the medical profession.\nRecently issued U.S. Pat. No. 4,224,946 describes a monofilament suture with good flexibility and knot strength, which suture is composed of segmented polyetheresters which contain (1) a polymeric block of polyalkylene ethers and (2) a polymeric block of aromatic dicarboxylic acids or cycloaliphatic acids with short chain aliphatic or cycloaliphatic diols. Similar subject matter is disclosed in Belgian Pat. No. 880,486.\nThe ethylene glycol polyester of the subject diacid moiety, 4,4'-(ethylenedioxy)bis benzoic acid, has been known for some time, (C.A. Registry No. [24980-45-8] if prepared from the acid, [26373-72-8] if prepared from the dimethyl ester) and, in fact, its ethylene glycol/polytetramethylene oxide copolymers (C.A. Registry Nos. [9071-04-9] and [51884-53-8] have been prepared:\nCA 76 11461Oy PA0 CA 80 83814u mentioned in index only PA0 CA 81 171028s PA0 CA 81 P12175g PA0 CA 83 P195002w manufacture of elastic fiber PA0 CA 84 P75556d polyester composite fibers with rubber-like elasticity) PA0 Knot strength--at least 20,000 psi PA0 Tensile strength--at least 30,000 psi PA0 Young's modulus--less than about 600,000 psi PA0 Elongation--from about 20% to 80% PA0 Knot strength--30,000 to 65,000 psi PA0 Tensile strength--45,000 to 100,000 psi (more preferably 60,000 to 100,000 psi) PA0 Young's modulus--75,000 to 600,000 psi (more preferably 75,000 to 250,000 psi) PA0 Elongation--from 20% to 55% PA0 a. burn dressings PA0 b. hernia patches PA0 c. medicated dressings PA0 d. fascial substitutes PA0 e. gauze, fabric, sheet, felt or sponge for liver hemostasis PA0 f. gauze bandages PA0 a. arterial graft or substitutes PA0 b. bandages for skin surfaces PA0 c. burn dressings (in combination with polymeric films) PA0 a. orthopedic pins, clamps, screws and plates PA0 b. clips PA0 c. staples PA0 d. hooks, buttons, and snaps PA0 e. bone substitutes (e.g., mandible prosthesis) PA0 f. needles PA0 g. intrauterine devices PA0 h. draining or testing tubes or capillaries PA0 i. surgical instruments PA0 j. vascular implants or supports PA0 k. vertebal discs PA0 l. Extracorporeal tubing for kidney and heart-lung machines PA0 m. artificial skin and others.\nas well as copolymers based on polyethylene oxide. These polyesters have been described as possessing improved crimpability, dyeability, and moisture absorption properties.\nThe 1,4-butanediol polyester of the subject diacid moiety is disclosed in Chemical Abstracts [52826-06-9]. In a 1978 paper (C.A. 89 60077c) which relates to the transesterification of dimethyl esters of aromatic dicarboxy acids with .alpha.-hydro-.omega.-hydroxy poly(oxyethylene)s and .alpha.,.omega.-alkanediols, a polytetramethylene oxide/1,4-butanediol copolymer based on the subject diacid moiety was made.\nIn spite of the fact that several polyethers could be incorporated chemically to toughen and lower the modulus of related polyesters, it is generally accepted that linear thermoplastics possessing a high initial modulus are more difficult to modify to increase compliance. Fibers of the (all hard) homopolymers of the subject invention were found to possess moduli in excess of 1.5 million psi, which renders it unsuitable for producing monofilament sutures.\nPatents that relate to (2-alkenyl or alkyl)succinates are U.S. Pat. No. 3,542,737 and U.S. Pat. No. 3,890,279. None of these patents discloses the present copolymers or modifications thereof.\nTheory and experience in the art of fiber chemistry predict that branching (such as that present in the instant copolymers) may inhibit fiber formation and will exert a deleterious effect on the tensile properties of any resulting fibers due to the inability of the unoriented branch to contribute to the load bearing capacity of the fiber; and by the stearic interference posed by the branch to chain alignment during fiber orientation. It is therefore surprising that strong fibers, in particular strong, flexible compliant fibers may be formed from the present copolymers with pendant hydrocarbon chains. As will be seen from Table 1, Example (i), homopolymers prepared from the 1,4-butanediol polyester of the subject diacid moiety, have a modulus of 1,678,000 psi. It is surprising that the modulus of certain of the copolymers of the present invention is reduced to only 64,000 psi.\nIt is an object of the present invention to provide a novel copolymer of poly(alkylene-4,4'-(ethylenedioxy)bis benzoate) as well as surgical devices formed therefrom. It is a further object of the present invention to provide a novel flexible, thermoplastic monofilament suture or ligature of said copolymer, having a diameter of from about 0.1 to 50 mil and possessing unique and desirable physical properties. Yet a further object of the present invention is to provide a Cobalt 60 sterilizable monofilament with lower modulus, better hand and more desirable tie-down characteristics than those of monofilaments of polypropylene. It is a further object of this invention to provide a monofilament suture with a desirable degree of ductility to accommodate changing wound conditions. It is yet another object of this invention to provide a monofilament suture with the flexibility and knot-tying characteristics of a braided suture. These and other objects will be made apparent from the ensuing description and claims."}
{"text": "Within the realm of social networking, the number of friends/followers/recipients is an indication of the relative level of success a user will have at sharing information. Famous individuals, corporate entities and other well established user accounts on social networking websites are constantly identified and have the potential to reach millions of friends/followers/recipients (for simplicity purposes the friends/followers/recipients will be referred to as data recipients or ‘recipients’). However, when the number of recipients reaches into the thousands, the social networking paradigm begins to break down.\nIt may be said that small and obscure groups of people who are passionately interested in a specific subject may be more likely to focus on finite tasks and accomplish their goals. Online social networking tools should be capable of fostering these sorts of clusters. But, when the conversation grows too large the focuses are often expanded and the recipients may lose interest quickly. Not only do audiences of recipients feel estranged, the participants also start self-censoring. People who suddenly find themselves with really large audiences often start writing more cautiously and behaving like politicians or corporate drones.\nIn one example, if someone or one user account has obtained 1.5 million recipients on a data messaging application, clearly they are among the rare and famous online individuals. Their response and data messages may include anything as they have already reached instant fame. In another example, if you have a hundred followers, you may be considered to just be merely chatting with pals. The middle ground—when someone amasses, roughly several thousand recipients, the social aspect of social media begins to fall apart by trying to keep their interests perked.\nFIG. 1 illustrates a graph depicting a number of data posts on a social networking user account in relation to a number of followers or data recipients of that particular user account. Referring to FIG. 1, it can be noted that the number of posts leveled-out after a user obtained 400 followers. According to the graph, it can therefore be theorized that once a user has more than 400 followers. The user account or sender did not increase the number of posts, but the number of posts remained around the 800 level.\nThe intimacy of the social network begins to change when the number of recipients grows into the thousands. For example, unlike when there are a limited number of followers, the user or sender begins to believe the recipients are satisfied with the content of the posts being sent. As the number of recipients grows, the user is hesitant to share the same intimate subject matter, but becomes naturally more political and concerned that the recipients may not be interested in the narrow scope of certain posts. As a result, the timeliness in receiving a response from a user's posts (feedback) may begin to diminish or disappear completely.\nA British anthropologist, Robin Dunbar, theorized that the maximum size of a social group in which everyone involved could maintain a mental record of all of the interpersonal relationships was 150. This is considered “Dunbar's number.” These are relationships in which an individual knows each person and how each person relates to every other person. Numbers larger than 150 generally require more restrictive rules, laws, and enforced norms to maintain a stable and cohesive group.\nDunbar's number is possibly one reason why users are more careful in posting messages to a large number of followers. Managing a large number of recipients is not easy, but the success and influence is directly proportional to the activeness of the recipients. Recognizing the behavior of the recipients may also provide insight into how to maintain their satisfaction with the content of the posts. Some of the factors that will determine the influential power of the posts and messaging may include the activeness of the recipients (feedback), the forwarding activity of the messages (sharing), the currency with which the recipients respond, etc.\nManaging a large and possibly increasing number of recipients may require certain modifications to the messaging application used to share information (e.g., social networking website, blog, group messaging, etc.). For example, with short messaging service (SMS) applications, the notion of groups includes followers or recipients who are tagged into different groups. The leader is then enabled to send data messages to one of the groups without the need of sending individual messages to each person, or sending a single message to all of the followers. Groups are viewed as a way to organize many people within a single group and to obtain more followers for each group. When a group is organized, all members of the group automatically follow the data messages posted for that group, and the leader may automatically follow the individuals as well. A leader may obtain more followers by joining a group. When a group is created, it is unique and no one can replicate that specific group. Therefore, the “administrator” of that group may be viewed as the expert. Anyone who joins that group will follow the leader or administrator of the group.\nGroups provide a social networking option for those leaders with a relatively small number of followers. The group situation becomes increasingly difficult to manage with a leader that already has a large number of recipients or is observing a significant increase in the number of recipients. Another concern is that the leader will not be able to divide all of the users into the correct group(s). The group strategy may be best utilized for leaders that do not have a large number of recipients.\nMultiple user accounts is another strategy to allow the leader to obtain a number of updates from a subset of recipients who are controlled by the leader. There are existing applications that allow leaders to maintain multiple messaging accounts and centrally manage those accounts. For example, when the leader is replying to the recipients, there is a function (e.g., a button on the “compose” message screen) that permits the leader to switch the “sender” or leader's account. The leader can setup a “follow”-only account, and when replying to something from that account, a simple sending option may permit the reply to be sent form the leader's “real” account. None of these examples provide the success desired by the social networking leaders and recipients. Many leaders are not satisfied with having multiple user accounts or managing each individual recipient among many others. The efficiency and productivity of managing a social network requires additional measures to be enacted."}
{"text": "The present invention pertains to a professional-type chair such as used by dentists, ophthalmologists, and other types of medical or quasi-medical professions in which adjustable operatory chairs are useful. Also, details and features of mechanisms found in such types of chairs are frequently employed in chairs used in barbering and beauty culture. One of the most common features now found in this type of chair comprises an upholstered seat to the rear edge of which a back rest is pivotally connected, and power means are employed to operate the back rest between substantially upright and rearwardly extending horizontal positions. Appropriate control means are employed to effect operation of one or more electric motors which effect movements of the back rest relative to the seat. In addition, many types of chairs of this type include a leg rest and/or a foot rest, the leg rest being pivotally connected to the seat adjacent the forward edge thereof and movement of the leg rest frequently occurs simultaneously with pivotal movement of the back rest relative to the seat.\nEffecting control of the back rest and/or leg rest relative to the seat is accomplished in certain modern types of chairs of the type referred to above by means of rather sophisticated electrical circuitry, batteries of push buttons by which various switches are operated to start and stop the movements of the back and/or leg rests relative to the seat, as well as raising and lowering the seat with respect to a supporting base. Examples of such control in chairs of this type are found in various patents, including U.S. Pat. No. 3,357,740 to Vaughn et al, dated Dec. 12, 1967, and U.S. Pat. No. 3,578,379 to Taylor, dated May 11, 1971. Notwithstanding such sophisticated control systems, especially for moving the back and leg rests relative to a seat, as disclosed in said patents, it now has been found that in adjustable operatory chairs of the type referred to above, it is desirable to provide adjustable, programmable control means by which, for example, any adjustable setting may be effected quickly and manually by the operator of the chair to dispose the back rest and/or leg rest at a desired position intermediately between the limits of movement thereof, such as between substantially vertical position and horizontally extending position, with respect to both the back rest and/or leg rest."}
{"text": "Cancer is the second most common cause of death in the United States, exceeded only by heart disease. In the United States, cancer accounts for 1 of every 4 deaths. The 5-year relative survival rate for all cancers patients diagnosed in 1996-2003 is 66%, up from 50% in 1975-1977 (Cancer Facts & Figures American Cancer Society: Atlanta, Ga. (2008)). This improvement in survival reflects progress in diagnosing at an earlier stage and improvements in treatment. Discovering highly effective anticancer agents with low toxicity is a primary goal of cancer research.\nProstate cancer is one of the most frequently diagnosed noncutaneous cancers among men in the US and is the second most common cause of cancer deaths with over 180,000 new cases and almost 29,000 deaths expected this year. Patients with advanced prostate cancer undergo androgen deprivation therapy (ADT), typically either by luteinizing hormone releasing hormone (LHRH) agonists or by bilateral orchiectomy. Androgen deprivation therapy not only reduces testosterone, but estrogen levels are also lower since estrogen is derived from the aromatization of testosterone, which levels are depleted by ADT. Androgen deprivation therapy-induced estrogen deficiency causes significant side effects which include hot flushes, gynecomastia and mastalgia, bone loss, decreases in bone quality and strength, osteoporosis and life-threatening fractures, adverse lipid changes and higher cardiovascular disease and myocardial infarction, and depression and other mood changes.\nMalignant melanoma is the most dangerous form of skin cancer, accounting for about 75% of skin cancer deaths. The incidence of melanoma is rising steadily in Western populations. The number of cases has doubled in the past 20 years. Around 160,000 new cases of melanoma are diagnosed worldwide each year, and it is more frequent in males and Caucasians. According to a WHO Report, about 48,000 melanoma-related deaths occur worldwide per year.\nCurrently there is no effective way to treat metastatic melanoma. It is highly resistant to current chemotherapy, radiotherapy, and immunotherapy. Metastatic melanoma has a very poor prognosis, with a median survival rate of 6 months and a 5-year survival rate of less than 5%. In the past 30 years, dacarbazine (DTIC) is the only FDA-approved drug for metastatic melanoma. However, it provides only less than 5% of complete remission in patients. In recent years, great efforts have been attempted in fighting metastatic melanoma. Neither combinations of DTIC with other chemotherapy drugs (e.g., cisplatin, vinblastine, and carmustine) nor adding interferon-α2b to DTIC have shown a survival advantage over DTIC treatment alone. Most recently, clinical trials with antibodies and vaccines to treat metastatic melanoma also failed to demonstrate satisfactory efficacy. Ipilimumab (Yervoy) is such a drug that uses your immune system to fight melanoma. Ipilimumab is used to treat advanced melanoma that has spread beyond its original location. Targeted therapy uses medications designed to target specific vulnerabilities in cancer cells. Vemurafenib (Zelboraf) is a targeted therapy approved to treat advanced melanoma that cannot be treated with surgery or melanoma that has spread through the body. Vemurafenib only treats melanoma that has a certain genetic mutation.\nTubulin/microtubule-interacting drugs are used successfully for treatment of a wide variety of human cancers. They are commonly classified into two major categories: microtubule-stabilizing (e.g., taxanes, epothilones) and microtubule-destabilizing drugs (e.g., vinca alkaloids, colchicine). Three major binding sites on α,β-tubulin subunits have been identified as taxane-, vinca alkaloid- and colchicine-binding sites. While antimitotic agents interacting with the taxane- or vinca alkaloid-binding sites in tubulin are tremendously successful in clinical oncology, there are no Food and Drug Administration (FDA)-approved colchicine-binding site drugs currently available for cancer treatment. Most of the colchicine-binding agents have high potency, relatively simple chemical structures for optimization, selective toxicity towards tumor vasculature, and show promising ability to overcome P-glycoprotein (P-gp) efflux pump mediated multidrug resistance.\nSeveral outstanding agents for such an approach are listed in FIG. 1. Combretastatin A-4 (CA-4) is the most active member of the combretastatins family, isolated from the African tree Combretum caffrum. CA-4 exhibits strong antitubulin activity by binding to the colchicine-site and has been the subject of Phase II and Phase III clinical studies. The replacement of the olefinic bridge of CA-4 with a carbonyl group yields phenstatin, which has similar potency and mechanisms of action with CA-4. BPR0L075 and Oxi-6196 are 2-aroylindole and dihydronaphthalene analogues of CA-4, which show strong inhibition on tubulin polymerization. Methylated chalcone SD400, which has an IC50 value of 0.21 nM against K562 human leukemia cells, is a potent tubulin inhibitor. Podophyllotoxin is a non-alkaloid toxin lignin, and it also possesses an anticancer property that can be attributed to the inhibition of tubulin polymerization through binding to the colchicine binding site. Accordingly, the colchicine-binding site compounds have attracted great interest from medicinal chemists in recent years."}
{"text": "The invention is based on a method for operating an electric power tool and an electric power tool according to the preambles to the independent claims.\nIn known electric power tools, a torque limitation is implemented in a known fashion in that either a mechanical overload clutch on the driven shaft of the transmission periodically interrupts the frictional engagement when a predetermined load moment is reached or the power consumption of the electric motor, which is proportional to the torque, is regulated to a constant value.\nAs a rule, the value of the desired load moment is preset and is not changed during operation. Usually, the user of the electric power tool makes changes when the motor is switched off. Through the actuation of a pushbutton of the electric power tool by the pressure of the user's finger, the speed of the motor is either changed by means of a pulse width modulation or is preset proportional to the pushbutton position by means of a regulator. When the load moment preset by means of the torque limitation is reached, the moment output by the motor is kept constant either by the mechanical clutch periodically interrupting the frictional engagement or by the motor current being regulated to a constant value.\nIn the case of the mechanical overload clutch, the motor continues to rotate at an undiminished speed; with the periodic engagement and renewed releasing of the overload clutch, a high torque peak is temporarily exerted on the work piece. In screwdrivers, this leads to the screw being turned a certain amount furthers e.g. with an impulse screwdriver, and can potentially lead to a stripping of the screw. The reaction of the overload clutch is also connected with a significant generation of noise that is often found to be unpleasant.\nIf the torque is electronically limited through regulation of the motor current, then after the preset load moment has been reached, the motor is supplied with current until the user releases the pushbutton. This also applies to the mechanical torque limitation."}
{"text": "This invention generally relates to a method and apparatus for controlling energizing interval of an ignition coil of an internal combustion engine, and more particularly, the present invention relates to such a method and apparatus for increasing the energizing interval during acceleration of the engine which is equipped with an electronic control device.\nGenerally speaking, conventional electronic control systems for internal combustion engines comprise various sensors, such as a rotational angle sensor, a reference position sensor, coolant temperature sensor, suction vacuum sensor, ambient air pressure sensor, acceleration sensor and the like, which detect the operating condition of the engine; a control unit for obtaining various amounts to be controlled, such as ignition timing, energizing interval of the ignition coil, fuel injecting interval, the amount of exhaust gas recirculation (EGR) and the like by using, for instance, a microcomputer, in accordance with the information from the sensors; and various actuators, such as an ignition device, an EGR control valve, a fuel injection device and the like, which control the operating condition of the engine in accordance with the control amounts from the control unit.\nIn the above-mentioned conventional electronic control system, the energizing interval of the ignition coil is determined by detecting the rotational speed of the engine crankshaft. Namely, the energizing interval decreases as the engine rpm increases. However, when the engine is accelerated, the actual energizing interval is apt to be shorter than a desired interval, as will be described hereinlater, in the conventional devices. Because of the shorter energizing interval, the ignition coil cannot produce a sufficient voltage for properly igniting the air/fuel mixture in the cylinders of the engine. Although it is theoretically possible to set a basic value of the energizing interval longer in advance so that the energizing interval during acceleration is sufficient, such a longer energizing interval may cause the final stage transistor and the ignition coil, both included in the ignition device, to generate undesirable heat, resulting in deterioration of these elements because an electric current flows through the transistor and the ignition coil for an interval longer than required during normal or steady operating condition of the engine."}
{"text": "Due to recent high interest of the 4th generation mobile communication system at home and abroad, developments of systems that satisfy the requirements thereof are being actively progressed. In particular, Orthogonal Frequency Division Multiplexing (OFDM) has been attracted by one of systems that are adapted to use in the 4th mobile communication system because of high transmission efficiency and simple channel equalization.\nAs typical multiple access schemes based on OFDM, there are OFDM-Frequency Division Multiple Access (OFDM-FDMA), OFDM-Code Division Multiple Access (OFDM-CDMA) and OFDM-Time Division Multiple Access (OFDM-TDMA). Among them, the OFDM-FDMA scheme allocates a different subchannel every user; and therefore, no Multiple Access Interference (MAI) is issued in intra cell and the number of subchannels and modulation level can be effectively changed depending on a transmission rate required by each user and channel circumstance. Thus; this scheme is advantageous to adaptive loading. In addition, since the OFDM-FDMA scheme is suitable for use of a great number of subcarriers, it can be applied to wireless communication systems with cells of wide area having relatively large time delay spread. And also, the OFDM-FDMA has drawn attention as a multiple access scheme for the 4th mobile communication system because it is suited to systems that support high vehicular speed.\nMeanwhile, duplex schemes are largely classified into Frequency Division Duplex (FDD) and Time Division Duplex (TDD).\nThe FDD scheme allows an uplink and a downlink to use different frequency bands wherein an interval between the frequency bands is enough as large as interference therebetween is negligible. This FDD scheme is mainly employed in macro cells that must provide relatively high time delay spread and high vehicular speed.\nThe TDD scheme shares a same frequency band for an uplink and a downlink, which is distinguished by setting a different time. This scheme can support services with asymmetric data transmission rate by modifying the number of slots of uplink and downlink and doesn't need to do a separate channel estimation of uplink and downlink by reciprocity characteristic of channel. However, such a scheme needs a guard time owing to round trip time and is mainly utilized in cells such as micro cells or hot-spot due to power problem.\nThe 4th mobile communication system is required to support high transmission rate and secure a variety of service qualities by using limited wireless resources. For this, there has been a need for developments of cell planning and resource allocation algorithms to increase cell capacity and efficiently allocate wireless resources to users.\nOn the other hand, new services are needed for mobility guarantee and high speed transmission under wireless environments. According to this need, technical development and standardization works for portable internet service of 2.3 GHz band and Mobile Broadband Wireless Access (MBWA) service are being progressed in the country.\nAs mentioned above, the OFDM scheme is one of the most spotlighted techniques due to the high transmission efficiency and simple channel equalization.\nMeanwhile, the most important property, which decides the performance of an OFDM multiple access-based system in cellular environments where mobility is considered, is Frequency Reuse Factor (FRF). This FRF implies the number of frequency sets used per cluster of mobile phone cells. In mobile phone business, it is important to accommodate the maximum subscribers with limited allocation frequencies, which is indirectly meant by FRF. As property of wireless scheme itself, an increase of FRF is possible by dividing service area into cells, micro-cells, pico-cells, etc. and again subdividing antenna sectors within a same cell, and with technique such as smart antenna, except for its own improvement.\nIf FRF is set as “1,” it may be most ideal in view of throughput of base station because it uses all wireless resources; but there May be serious performance degradation due to inter-cell interference caused by FRF “1.”\nTo implement FRF “1” while solving the performance degradation problem by the inter-cell interference, there is developed a flash-OFDM system by the Flarion company. This flash-OFDM system employs a scheme that maximally prevents the performance degradation by the inter-cell interference with Low Density Parity Check (LDPC) channel code by using a frequency hopping scheme that changes subcarriers of OFDM with a constant pattern. In addition, there is a scheme that randomly punches subcarriers to decrease collision of themselves and their adjacent cells and so as to implement FRF “1.”\nIn the systems that maintain FRF as “1”, however, the performance degradation still exists at the boundary of cells where channel condition is poor due to the inter-cell interference by increase of traffic load. Accordingly, as method for securing the performance of user terminals in an area where channel condition is poor such as the boundary of cells, together with improvement of frequency efficiency, there is a frequency reuse partitioning scheme.\nThe frequency reuse partitioning scheme is one of effective methods capable of elevating the frequency efficiency.\nFIG. 1 is a view of describing a general frequency reuse partitioning scheme.\nBasic idea of the frequency reuse partitioning scheme is that cells 111 to 117 are divided into an inner cell 101 and an outer cell 102 based on the length between a base station and a user terminal or an intensity of pilot signal sent from the base station to the user terminal, and then different FRFs are applied to the inner cell 101 and the outer cell 102.\nAs may be seen from FIG. 1, if the user terminal exists in the outer cell 102 area, subchannels with FRF “7” are allocated to cells 111, 112, 113, 116; and if the user terminal is within the inner cell 101 area, subchannels with FRF “1” are assigned to cells 114, 115, 117.\nAs mentioned above, the reason the channels with different FRFs are allocated to the inner cell 101 and the outer cell 102 is because it may be an importance factor that restricts cell radius of cellular system for the following: in case of terminals near to a base station, channel state is good since power loss by path loss is less than that of terminals far from the base station, while, in case of terminals around the boundary of cells, performance degradation is occurred and data transmission rate is constrained due to serious affection by power loss by path loss and inter-cell interference.\nTherefore, in use of the frequency reuse partitioning scheme, channels with low FRF are assigned to the user terminals 114, 115, 117 in the inner cell 101 area to secure a proper level of service quality since their channel state is generally good, thereby enlarging the capacity of cell. Meanwhile, channels with high FRF are allocated to the user terminals 111, 112, 113, 116 in the outer cell 102 area because their channel state is relatively poor so that a radius of cell can be expanded and thus the same level of service quality and data transfer rate as that of the terminals in the inner cell can be secured, in case of the user terminals located at the boundary of cells.\nIn addition, wireless resource allocation methods have been recently studied to effectively employ limited frequency resources while alleviating inter-cell interference.\nIf it is assumed that channel is stationary and channel response of user is accurately perceived by transmission end, it is known that a combined scheme of water filtering and adaptive modulation technique is optimal one. This water filtering technique has been mainly developed only in single-user systems or multi-user systems that support constant resource allocation: for example, a system using TDMA or FDMA allocates slots or frequency channels for a given time for each user and then applies an adaptive modulation technique for channels of each user.\nIn the multi-user OFDM scheme that adopts the adaptive modulation technique based on the constant resource allocation as described above, however, there is a drawback in that it cannot know optimal resource allocation that can be provided by actual system. The reason is because there are many unused channels when water filling algorithm is applied due to presence of subchannels that suffers from deep fading or subchannels to which high power cannot be given in view of frequency selective channel characteristic.\nHowever, there may be an instance where channel appeared as deep fading to one user may not be appeared as deep fading to another user. Generally, if the number of users is increased, a probability that each of subchannels constituting OFDM is appeared as deep fading channel to all users is gradually decreased. In other words, a multi-user diversity gain can be acquired by bearing independent channel as the number of users is increased.\nConsequently, there has been required a study of an improved scheme capable of enabling an optimal resource allocation by dynamically allocating relatively good channel to each user based on channel information of all users and then applying an adaptive modulation scheme using those channels."}
{"text": "1. Field of the Invention\nThe present invention relates generally to connectors which can be mated and unmated in very harsh environments, such as underwater, and is particularly concerned with a robotically mateable, rotary joint electrical connector usable in such environments.\n2. Related Art\nThere are many known electrical rotary joint connectors which function in dry environments, but not many that are suitable for harsh or underwater environments. Harsh environment rotary joint connectors are currently in use only for limited applications, such as down-hole drill strings, but they are not intended for general use, or for mating while completely submerged in a hostile environment.\nModern underwater systems, particularly sea floor systems, are generally of modular architecture with individual modules added and removed in-situ, and these modules are typically electrically interconnected by various jumper cables which have underwater mateable and demateable connectors on one or both ends. The connectors consist of two mating halves or units, specifically a plug unit and a receptacle unit. When these connector units are brought together in the mating process, the following relative alignment aspects must be controlled: axial offset (along the mating axis), axial tilt (of the mating axes), and rotation about the mating axis. The last of these, which requires rotational keying, is the most difficult to control.\nAll jumper cables have some resistance to torque, so that if one end is fixed, or substantially fixed due to great length, it is often very difficult to rotationally orient the free end into mating position, particularly if using the robotic manipulator of robot or robotic vehicle, such as an underwater vehicle or remotely operated vehicle (ROV). The problem is increased for very short jumper cables which cannot support much twist. Even in conditions where there is little cable torque, the problem of finding the mating key/key slot position increases the difficulty of connecting the jumpers. It would therefore be advantageous to provide a system where no keying of the mating connector portions or units is required, and which avoids the problem of twisting of the jumper cable when the connector on the cable's free end is rotated, as happens due to underwater vehicle maneuvering.\nMost current underwater connectors designed for mating while completely submerged need to have the plug and receptacle parts keyed into rotational alignment for mating. Many of the connectors used for the oil and gas industry's subsea operations are connected and disconnected remotely, either by being mounted to large, opposed plates (stab plates) that come together during the mating process to join arrays of connectors, hydraulic couplers, and the like, or by robotic manipulators such as manipulators of remotely operated vehicles (ROV's). Mating remotely is made more difficult and expensive by the requirement to control the rotational alignment of the individual components to be mated.\nIn the early 1980's two-contact fluid-filled electrical connectors that required no rotational alignment were made commercially available. One example of such a connector is described in U.S. Pat. No. 4,606,603 of Cairns. These connectors did solve the rotational alignment problem, but one problem with such connectors was that two contacts were not enough to satisfy the needs of most operations. Another problem was that the receptacle's circular end-opening, which had to be pinched tightly closed before and after mating, had to be stretched several hundred percent to receive the plug's pin. If mated for a long time, particularly at low ocean temperatures, the opening did not close upon de-mating, and the connector subsequently failed.\nIn the late 1980's, multiple pin, fluid-filled connectors were once again introduced. They have all the required barriers, are robust, and exceptionally reliable. One such connector is the subject of U.S. Pat. No. 4,948,377 of Cairns. These connectors are manufactured by Teledyne ODI. They replaced the two-contact, single pin fluid-filled connectors described above as the high-reliability standard for the offshore industry. These connectors still have the rotational alignment problem, however, which somewhat limits their use, and requires special keying provisions for rotational alignment.\nIn the early 1990's a keyless, coaxial, oil-filled, wet-mateable connector was introduced that required no rotational alignment. This connector is described in U.S. Pat. No. 5,171,158 of Cairns (hereinafter '158 patent). It consisted of multiple ring-like contacts spaced along a constant diameter portion of the plug pin. The receptacle had corresponding ring-like contacts spaced along a rubber bore to receive the plug contacts. The overall layout of the contacts was very similar to the first type of connector described above. The main differences were that the connector of the '158 patent housed the receptacle contacts in a pressure-balanced, fluid-filled chamber; and, when mated, the individual pin/socket pairs were separated from each other by a single rubber seal. Unlike the coaxial connector of U.S. Pat. No. 4,606,603 (hereinafter '603 patent), the anterior sealed opening through which the plug's probe passed when entering the receptacle's chamber was occupied by a spring loaded piston before and after mating. That removed the necessity of the sealed opening to be pinched closed to a zero diameter as in the '603 patent.\nThe connector shown in the '158 patent was reasonably successful technically and quickly gained a dedicated customer base, but it was discontinued after being on the market for just a couple of years. It proved to be too expensive and difficult to manufacture. It also still had the problem of cross-connection during mating and de-mating as the plug's contacts wiped across receptacle contacts which were not their intended counterparts.\nUnderwater and other harsh environment connectors today typically require rotational alignment or keying for connection in a harsh environment, which makes robotic mating and demating difficult to accomplish, and also often results in torque or twisting of the flying lead jumper cables."}
{"text": "Conventional simulation methods to solve structural dynamics problems are in the domain of finite element technology where the problem is solved in the modal domain and then the results are mapped into the time domain by appropriate transformations.\nLimitations of the prior art have been in the difficulty in deriving new three-dimensional elements for different applications of interest and the lack of ease in obtaining the temporal solution directly from the solution of governing elastodynamic pdes.\nThe need to know the state of a structural system during its operation in terms of the physical output variables such as stresses and the geometric configuration of the system itself is essential for monitoring the system health. Such systems can be tested, prior to launching them in their operational domain, in a laboratory or through relatively inexpensive computational simulations.\nSuch systems when subjected to space and time varying loads during their operation can throw the system into unsafe states from the system's health perspective. It is therefore essential to have a prior knowledge of such system states before the systems are commissioned.\nFor the reasons stated above, and for other reasons stated below which will become apparent to those skilled in the art upon reading and understanding the present specification, there is a need in the art for alternative modeling and simulation methods for structural systems."}
{"text": "Today, companies and persons rely on having consistent supply of power to electronic devices more than ever. Without power, companies may be unable to manufacture goods, or to operate at all, such as if the company is in the business of supplying information over the Internet. Without power, businesses and individuals may be completely incapacitated regarding critical activities, such as making goods, providing services, and transacting personal finances (e.g., filing tax returns, and paying bills).\nFrequently, individuals and companies monitor the supply of power to electronic devices or facilities to ensure that their affairs and/or businesses are not significantly affected by a power outage, a change in status or some other power alteration. Uninterruptible power supplies (UPSs) are often used to provide backup power in case of a power outage. UPSs are commonly used on computing equipment to guard against data being lost due to a power outage before the data are saved. UPSs used with computing equipment also help to guard against a loss in service by providers of information over the Internet, such as by servers, e.g., hosting web pages.\nDevice status is often monitored and reported to the individual or company owning a device using email notification or the Plain Old Telephone System (POTS), which requires the installation of a remote monitoring unit at the customer site and is physically dialed into to get updates and specifics. Frequently, companies are obligated to provide manpower at substantial cost so that electronic devices can be monitored and reported on a regular basis."}
{"text": "In the past, many corn-based products were prepared using continuous deep-fat frying techniques. However, products cooked using oil contained cholesterol or fat and are generally deemed to be less wholesome than those without such by-products. To remove such undesirable contents, the industry attempted to produce the same type of product without using oils for cooking by baking the product. However, by changing the cooking processes, the taste of the product also changed.\nMany factors play a role in the way a product tastes. These factors include ingredients, moisture content, and cooking time. The same factors also affect the texture of the product such as crispness, crunch, or bite, and the like. No matter whether the product was cooked by deep-fat frying or by baking, the industry has always thoroughly mixed the ingredients to produce a homogenous dough. To produce the dough, the ingredients were mixed sufficiently prior to sheeting so that the resulting dough was readily workable and would form an even sheet.\nFully gelatinizing the dough by thoroughly mixing and cooking the ingredients produces a glassy sheet of dough having a high moisture content, which requires longer baking and drying periods to drive off the moisture. Typically, such results produce a product which is brittle and has a less-desirable texture and taste."}
{"text": "For some years there has been a move for limiting and restricting smoking in certain public areas. For example, there has long existed smoking and non-smoking areas on commercial airlines. Moreover, the Civil Aeronautics Board is now considering totally banning smoking on short flights. Smoking restrictions such as those presently found on commercial airlines exist in other public areas as well. Many State and local governments have considered and are presently considering legislation aimed at restricting smoking in certain public areas, especially public areas where individuals find themselves in close quarters. Thus in the future it is likely that such restrictions on smoking will have a substantial effect on smoking and the right and ability to smoke in public areas.\nWhile restricting smoking in public areas may be effective, it is interesting to consider that very little effort has been made in the area of designing a smoking apparatus or a smoking system that will enable an individual to smoke tobacco material or a cigarette and not bother individuals in and around the individual smoking. In this regard, what is being alluded to is a smoking apparatus that can be deemed to be a closed or generally closed smoking system. Such a closed smoking system entails a smoking apparatus wherein the smoking material or cigarette contained within a structure and during the smoking process the resulting smoke from the cigarette and/or the individual smoking is somehow confined.\nThere has been some very limited efforts at designing such a generally closed smoking system. For example, see the disclosures found in U. S. Pat. Nos. 4,198,992; 4,211,244; 4,198,992 and 1,792,279. However, these smoking devices are not truly totally closed smoking systems. In addition they have tended to be large, bulky and hard to handle. In addition besides being large and bulky, such smoking devices of the prior art have tended to be complicated and hence expensive."}
{"text": "The present invention concerns a method of machining rotationally symmetrical parts to put them into the condition in which they are capable of use, wherein the parts in the course of the manufacturing procedure are at least partially subjected to a hardening procedure. In particulars this concerns crankshafts, in particular the bearing surfaces (both of the big-end bearings and also the fitting bearings, the sealing bearings and central bearings) of crankshafts, in regard to which the condition in which they are capable of use is that condition in which the crankshaft can be fitted in the engine without further removal of material at the bearing surfaces. In that respect bearing surfaces is used to denote both the peripheral surfaces, that is to say the bearing diameter, and also the surfaces which are referred to as the mirror or thrust surfaces, that is to say the end faces which adjoin the peripheral surfaces generally in perpendicular relationship and which serve for example for axial thrust support. In this respects particularly in the case of steel crankshafts, the bearing surfaces are hardened in the regions thereof which are in the proximity of their surface. That serves both for increased resistance to abrasive wear of the bearing locations, and also serves to protect those bearing surfaces from damage in the course of handling throughout the entire manufacturing procedure as well as serving to influence the strength properties involved.\nCrankshafts, in particular the crankshafts for private motor vehicles which have a large number of cylinders, are known to be workpieces which are unstable during machining and thus difficult to machine. Assessment of the dimensional accuracy of a finished crankshaft is effected primarily, besides the axial bearing width, by assessment of the following parameters:\nDiameter deviation, which is deviation from the predetermined reference or target diameter of the bearing journal,\nroundness, which is macroscopic deviation from the round reference or target contour of the bearing journal,\nconcentricity, which is a diameter deviation in the case of a rotating workpiece, that is to say, for example, the deviation from the reference positional contour which a main bearing location effects during the rotary movement of the crankshaft by virtues on the one hands of the non-ideal roundness of that main bearing journal and, on the other hand, of flexing of the crankshaft which in that case is clamped only at its ends,\nroughness Ra, which is a value which is ascertained by calculation and which represents the microscopic roughness of the surface of the bearing location, and\npercentage contact area, which is the load-bearing surface proportion of the surface structure, considered microscopically, which comes into contact with a co-operating or counterpart surface pressed thereon.\nIn addition, in regard to the big-end bearing locations, there is an assessment of:\nstroke deviation, which is the dimensional, percentage deviation of the actual stroke (spacing of the actual center of the big-end bearing journal from the actual center of the main bearings), from the reference or target stroke, and\nangle deviation, which is the deviation of the actual angular position of the big-end bearing journal from its reference or target angular position relative to the main bearing axis and with respect to the angular position in relation to the other big-end bearing journals, the angle deviation being specified in degrees or as a length dimension in the peripheral direction, related to the stroke.\nObserving the desired tolerances in regard to those parameters is made difficult less due to the available machining methods than the instability of the workpiece and the machining forces involved. The efficiency and economy of the method also play a large part in a practical context.\nHitherto the removal of material from the bearing locations on the crankshaft in its original form, that is to say as cast or forged, was effected in succession in four machining steps:\nFirst Step:\nCutting machining with a given cutting edge; this involved using the processes of turning, rotary broaching, turning-rotary broaching, internal round milling and external milling, rotary milling, in particular in the form of high-speed milling or combinations of such procedures. The magnitude of the material to be removed was in the millimeter range.\nSecond Step:\nHardening of the desired regions, in particular the bearing surfaces, wherein heating and subsequent cooling of the crankshaft and the change in the material and configurational structure, in particular due to the irregular way in which those procedures take place, causes distortion and dimensional variation of the crankshaft or one of the regions thereof, and that has to be compensated by the subsequent operations of removing material, in consideration of the reference or target dimensions of the final condition.\nThird Step:\nGrinding by means of a hard, massive grinding tool, for example a grinding wheel, which generally rotates with its axis of rotation in parallel relationship with the axis of rotation of the crankshaft to be machined; the amount of material to be removed was in the tenths-of-millimeter range.\nIn the case of crankshafts which are difficult to machine, in particular crankshafts which are long and thus highly unstable, the grinding machining operation was also effected in a multi-stage procedure, for example in a two-stage procedure by preliminary and finishing grinding.\nFourth Step:\nFinishing by generally a stationary grinding means (grinding belt or grinding stone) which is pressed against the external periphery of the rotating bearing location; the amount of material to be removed is at the present time in the range of hundredths of a millimeter or even mm.\nIn that respect, a distinction is also to be drawn in regard to the machining operation, in respect of the material of the crankshaft (steel or cast iron), in which connection in particular steel crankshafts which are preferably used for situations of use involving a high loading are hardened after the cutting machining operation.\nIn order to reduce the costs involved in crankshaft machining, the endeavour is to reduce the machining of the bearing locations from four to three different machining stages.\nThis means however when the cutting machining operation is omitted that in particular the removal of material which is to be implemented by the grinding operation must be greater than in the case of the four-stage method.\nRemoving material by means of grinding however involves the following disadvantages:\nbecause of the cooling/lubricating agent which is to be added, the grinding slurry which is produced gives rise to problems and is extremely costly to dispose of,\nbecause of the oil contained in the cooling/lubricating agent, for example in the case of CBN-grinding, there is always a latent risk of explosion,\nin the grinding operation the amount of cooling/lubricating agent used is substantially greater than in the case of cutting machining procedures as the cooling/lubricating agent is additionally employed in order to remove the grinding dust and swarf from the surface of the grinding wheel again, by jetting the cooling/lubricating agent on to same under high pressure,\nnonetheless the danger of the workpiece suffering from overheating is very high,\nthe machining pressures acting on the workpiece are higher than in the case of cutting machining, and\na microscopic surface structure is produced, in which the grain boundaries which are torn open by the grinding grain are smeared closed again by the subsequent grinding grains, with removed workpiece material, that is to say this is a surface structure with relatively few steep peaks, but with more or less flat, bent-over peaks which partly overlap in scale-like relationship.\nA primary object of the invention is to simplify the removal of material when machining bearing locations on a crankshaft to be hardened.\nThis object is attained by virtue of the grinding machining operation being omitted so that the machining sequence is reduced from four to only three machining procedures which are, in principles different. This eliminates not only all disposal problems in regard to the grinding slurry or swarf, but also the quite considerable capital investment costs for grinding machines, the costs involved in tool consumption and not least the required stock of workpieces, which is increased due to the grinding operation, by virtue of prolonged turn-around times for the workpieces. Disposal of the cuttings or swarf from the cutting machining operation does not give rise to any problems as either cutting is effected dry (high-speed milling) or separation of oil and cuttings or swarf is entirely possible by virtue of the much lower specific surface area of the cuttings or swarf in relation to grinding dust.\nSo that, in the procedure for removing material, the finishing operation can directly follow the cutting machining operation with a given cutting edge, hereinafter referred to for the sake of brevity as cutting machining, the degree of the admissible deviation of the actual values from the reference or target values, as occur after the cutting machining operation, must be so established that, in the totality of the machining procedures (cutting machining+finishing), the complication and expenditure involved must be technically as low as possible, with at the same time an overall machining time that is as short as possible.\nUnder some circumstances it is also possible to effect cutting machining with such a good final result that it is possible completely to omit the finishing operation. Overall the present invention therefore concerns the following combinations of processing sequences until the condition in which, for example, the crankshaft is capable of use is achieved:\ncutting-hardening-cutting-finishing\ncutting-hardening-cutting\ncutting-hardening-finishing\ncutting-hardening.\nIn that respect it is not sufficient in the cutting machining operation to strive for reference or target dimensions which come as close as possible to the final dimensions after the finishing operation, so that the oversizes which are to be dealt with by finishing and thus by relatively slow removal of material can remain as small as possible.\nIt must be taken into account that the stroke deviation and the angle deviation of the big-end bearing journals can no longer or can be only very slightly compensated by the finishing operation.\nIt must also be borne in mind that in the finishing operation in any case firstly the amounts of material that can be removed (reduction in diameter) are very small, that is to say up to about 200 mm can be achieved at economically viable expense and in addition secondly the finishing operation primarily provides for an increase in the percentage contact area, more specifically by a reduction in roughness, with the aim of achieving a percentage contact area of about 95%. A percentage contact area of 100% is unwanted as then there would no longer be any depressions which are necessary in order to maintain a film of lubricant at the bearing.\nIn regard to the microscopic surface structure, the cutting machining operation gives a surface in which the grain boundaries are partially torn open by virtue of the cutting edge pulling the grains apart as it cuts through the material. As a result the surface has a relatively large number of pointed raised portions, interrupted by valleys in the form of opened grain boundaries. A surface structure of that kind is conducive to the finishing operation by virtue of the fact that the many pointed raised portions not only facilitate the removal of material by the finishing procedure, but at the same time they also delay clogging of the finishing belt and the like member which is used in the finishing operation, insofar as the pointed raised portions of the workpiece surface provide that the material which is already deposited in the finishing belt is partially torn out of same again.\nThe limitations of this procedure are already encountered in the operation of reducing roundness deviations by finishing, insofar as the amount of time involved or the roundness deviations which can be dealt with depend not only on the absolute value of the roundness deviation to be equalised or levelled, but also the configuration thereof:\nIf the non-roundness is such that there are only a few (for example 2-3) troughs and raised portions, distributed over the periphery (thus constituting long-wave non-roundness), then, with the same absolute value in terms of non-roundness, for equalisation by means of finishing, a substantially greater amount of time is required or, under some circumstances, it is not possible to provide for complete equalisation, in comparison with short-wave non-roundness involving about 5 and preferably even about 30 or more troughs per periphery of the bearing location, with the same absolute value in respect of non-roundness.\nIt is also to be borne in mind that in the finishing operation at the same time on the one hand the degree of roughness is reduced and thus the percentage contact area is improved, while on the other hand the existing non-roundness is equalised or levelled. Those two effects can scarcely be decoupled from each other, or can be decoupled only to a very limited extent. If therefore, starting from an initial roughness, the desired roughness is achieved in the finishing operation after a given period of time, the finishing procedure is stopped as a given percentage contact area should not be exceeded. The equalisation or levelling effect in respect of the roundness deviation, which is achieved in that condition, is then accepted as a final result, and cannot be advanced separately any further.\nAccordingly, if the procedure is commenced from a given initial condition in respect of those two parameters from the finishing operation, roughness and percentage contact area cannot be machined independently of each other to afford desired final values.\nFor the direct succession of the finishing operation after the cutting machining operation therefore what is recommended is in particular specific coupling of the input parameters in regard to the finishing operation and therewith the output parameters in regard to the cutting machining operation, in respect on the one hand of the absolute value and degree of roundness deviation and on the other hand in respect of microscopic roughness and the percentage contact area which applies in that case.\nIn the cutting machining procedure using turning methods and also rotary broaching methods and possibly also when using slow milling methods, long-wave roundness deviations rather occur by virtue of long-wave oscillations in the machine structure, the tools and the workpieces involved in cutting machining. In contrast more especially high-speed external milling gives rise to short-wave roundness deviations. Therefore, the use of high-speed external milling in which an external round milling tool of a diameter of about 700 cm, which is very large in diameter in comparison with the crankshaft, rotates at a cutting speed of between 150 and 1000 m/min beside the relatively slowly rotating workpiece about an axis which is parallel with respect to the crankshaft means that it is possible to achieve roundness deviations with many raised portions along a circumference of a bearing.\nRotary milling, in particular if it is implemented with high cutting speeds, also rather tends to involve short-wave oscillations and thus short-wave roundness deviations. For, rotary milling involves milling with a kind of end-milling cutter which is arranged in parallel displaced relationship with respect to the radial direction of the peripheral surface, which is to be machined, of the bearing location, insofar as the peripheral surface is machined by means of the one or some cutting edges arranged on the end of the end-milling cutter. In that case, in particular machining with only one single cutting edge has been found to be advantageous if in that case operation is implemented with very high speeds of milling cutter rotation and the workpiece rotates comparatively slowly. If in that case the above-mentioned mirror surfaces or thrust surfaces of a bearing location are also to be subjected to machining, the end-milling cutter is also provided with one or more cutting edges on its peripheral surface.\nIt must also be taken into account that, in the finishing operation, the grinding means which bears against the workpiece, for example a finishing belt, is generally not changed during the procedure. The grinding means therefore becomes increasingly clogged at its surface during the finishing operation and the amount of material removed per unit of time progressively decreases.\nHow fast the grinding means begins to suffer from clogging in particular at the beginning of the finishing operation depends not only on the initial roughness of the surface but also on the percentage contact area thereof:\nThe lower the percentage contact area with a given level of roughness at the beginning of the finishing operation, that is to say the more pointed the microscopic surface structure is with correspondingly steeper flanks, then all the more readily can the material particles which have been removed from the surface and deposited in the finishing belt be removed from the finishing belt or the finishing stone or comparable finishing means at the beginning, when dealing with such a surface structure. With increasing equalisation or levelling of the microscopic surface the surface of the grinding means also becomes progressively more quickly clogged in the finishing operation.\nThis means that, with the same degree of roughness, a percentage contact area which is low at the beginning of the finishing operation is advantageous for high initial removal of material and thus a finishing operation which is short in time.\nThis also means that the degrees of roughness which can be handled by the finishing operation increase, in inverse proportion to the percentage contact area involved with those greater levels of roughness.\nWith the previous grinding operation, which was always effected following hardening, the roughness of the surface was admittedly reduced in comparison with preliminary cutting machining, but in-that respect in the same way the percentage contact area was either kept constant or even increased as the preliminary cutting machining operation left behind a microscopic surface structure which involves a low percentage contact area as machining with the given cutting edge, in the regions near the surface, means that in part the grain boundaries in the metal structure are torn open, extending radially from the outside inwardly.\nIn that way it is possible to finish in an economic fashion directly after the cutting machining operation insofar as efficient removal of material in the finishing operation is promoted and assisted by the preliminary machining procedure, insofar as on the one hand the choice of the correct cutting machining procedure means that the roughness achieved in that case has a low percentage contact area, and the roundness deviation achieved in that case is a roundness deviation which is as short-wave as possible.\nIn that respect it must further be taken into consideration that, in the previous grinding of bearing locations, the roundness deviations resulting from the preliminary cutting machining were generally only reduced by the grinding operation in terms of their absolute value but not in terms of their characteristic. Therefore, the grinding operation did not result in long-wave roundness deviations becoming short-wave roundness deviations, but the number of troughs was either retained or even reduced, with the consequence that a further improvement in roundness deviations by the finishing operation, considered as an improvement in result per unit of time, was made rather more difficult in the finishing procedure.\nThat means that a finishing operation directly after the cutting machining procedure is particularly economical when, after the cutting machining operation, the roundness deviations are less than 60 mm and in particular less than 40 mm, the diameter deviation is less than 200 mm, in particular less than 150 mm and the roughness Ra is less than 10 mm, in particular less than 6 mm. In that respect, the aim to be sought is roundness deviations with a short-wave nature of at least 5 waves per circumference, which applies in regard to bearing diameters of about 50 mm, but which, with rising or falling bearing diameters, should only change in a sub-proportional fashion, that is to say for example a 100% change in circumference produces only a maximum 30% change in the number of waves.\nFurthermore, in that respect, the aim to be sought is a rather lower percentage contact area in respect of the roughness achieved after the cutting procedure, than is obtained after the grinding procedure.\nIn the case of big-end bearings in addition the angle deviation after the cutting machining operation should be less than 0.4xc2x0, in particular less than 0.2xc2x0, and the stroke deviation should be less than 0.40%, in particular less than 0.20%, which corresponds to the tolerances to be observed in regard to the crankshaft when ready for use, as those parameters can no longer be changed in an economically viable manner by the finishing operation.\nA suitable form of the cutting machining procedure is therefore external milling or rotary milling, in particular in the form of high-speed milling, in consideration of the above-described interrelationships.\nParticularly when dealing with heavy workpieces, it can be found that a combination of the deviations between reference or target values and actual values of the relevant parameters, such deviations occurring after the cutting machining operation and prior to the finishing operation, wherein such combination is desirable in terms of direct coupling of cutting machining and finishing, can only be achieved if the cutting machining procedure is effected in a plurality of stages, in particular in two stages (preliminary cutting and finishing cutting). In that respect, high-speed external milling or high-speed rotary milling is to be preferred both when dealing with big-end bearings and also when dealing with main bearings, for the second stage, the finishing cutting operation.\nThe first stage, the preliminary cutting operation, when dealing with big-end bearings will also be effected using external milling, in particular using high-speed external milling, while when dealing with main bearings this can also be effected by turning or rotary broaching or turning-rotary broaching.\nIf the cutting operation is implemented in two or even more stages, the oversizes which are to be dealt with in the finishing cutting operation range in the optimum fashion between 0.1 and 0.4 mm in order further to improve in particular roundness and diameter deviation by virtue of the removal of very thin cuttings or swarf and in order to achieve a roughness which remains as equal as possible in the course of each individual step in the high-speed milling procedure, with the percentage contact area remaining uniformly low, from a microscopic point of view.\nWhether hardening is effected only after conclusion of the single-stage or two-stage cutting machining procedure, and material can be removed after the hardening operation only by means of finishing, depends primarily on the extent of the distortion which occurs in the hardening operation. The distortion depends on a series of individual factors such as configuration and dimensioning of the crankshaft, the nature of the hardening process, whether relatively small or relatively large surface regions or indeed the entire surface are or is hardened, the depth of penetration of the hardening procedure and the configuration in respect of time of the hardening procedure, and so forth. In that respect the probable maximum distortion which occurs in the hardening operation, together with that after the last cutting machining step effected prior to the hardening operation, is not to exceed the threshold which can be managed at commercially acceptable cost by finishing.\nOn the other hand after the hardening operation a material-removal step is again implemented by cutting machining, in particular the second cutting machining step (finishing cutting) of a two-stage cutting procedure.\nAs the removal of material in the cutting machining operation is more difficult after hardening than when working on the unhardened workpiece, in this case the oversize which is to be dealt with by the finishing cutting operation in the hardened condition is kept at a small level.\nThe use of belt-type finishers is recommended for the finishing operation, in which case the grinding belts are pressed by means of contact pressure shell members against the rotating bearing location and at the same time a relative oscillation is produced between the grinding belt and the workpiece, in the longitudinal direction. In that case, the contact pressure shell members should embrace the workpiece by at least 120xc2x0 in each case, preferably by up to 180xc2x0, of cooling/lubricating agents to the machining location, although this cannot always be achieved."}
{"text": "The present invention relates generally to motivational development of desirable behavior. More particularly, the present invention relates to a motivational apparatus which allows use of an electronic device only after a predetermined number of tasks have been successfully completed.\nAttempts have been made in society to use motivational methods to achieve certain goals. Much of human behavior can be understood as being directed toward specific goals. Incentive motivation is concerned with the way goals influence behavior. One of the most important aspects of this type of motivation is that any goal one seeks can motivate behavior. For example, the goal of obtaining spending money can serve as a strong motivator for performing yard work.\nCommonly, parents will require their children to perform tasks or xe2x80x9cchoresxe2x80x9d in order to obtain privileges. The child\"\"s ability to obtain a privilege is typically made contingent on the performance of socially desirable acts, such as cleaning up one\"\"s bedroom, good personal hygiene, being obedient, and being generally respective of others.\nIn today\"\"s society, many of the privileges enjoyed by children are related to electronic devices, such as watching television, playing computer games, and talking on the telephone with friends. Making these privileges contingent on the performance of certain tasks would provide a strong motivation to perform these tasks, which, over a period of time, can develop socially desirable behavior and an appreciation for delayed gratification.\nAccordingly, what is needed is an apparatus which promotes development of desirable behavior by providing a motivational apparatus which allows access to electronic devices only after a predetermined number of tasks have been successfully completed.\nIt is an object of the present invention to provide an apparatus for controlling use of electronic devices.\nIt is another object of the present invention to provide a motivational apparatus which allows use of an electronic device only after a predetermined number of tasks have been successfully completed.\nIt is another object of the present invention to provide a method for using the motivational apparatus.\nThese and other objects of the invention are accomplished with a motivational apparatus for controlling use of an electronic device, wherein the apparatus comprises means for interconnecting the apparatus between a source of transmission and the electronic device, and means for controlling transfer of transmissions from the transmission source to the electronic device. The transmissions may be electrical power from an AC or DC source, radio or video signals, analog or digital signals, or any other form of energy or information. The controlling means preferably comprises a plurality of electrical switches, a master switch, and a plurality of light-emitting diodes. Each of the electrical switches is preferably operable to open and close an electrical circuit to activate a diode to emit light and to open and close a portion of a master electrical circuit. The master switch, actuation of which is preferably restricted, is operable to open and close a portion of the master electrical circuit. The electrical switches must all be actuated to allow transfer of transmissions from the transmission source to the master switch, which, in turn, must be actuated to close the master circuit and thereby allow transfer of transmissions from the transmission source to the electronic device.\nIn operation, the apparatus is interconnected between the electronic device and the transmission source. A list of tasks is prepared and associated with the electrical switches such that one task corresponds to one electrical switch. As each task is successfully completed, the electrical switch corresponding to that task is actuated to light up the diode corresponding to that switch. After all of the tasks are completed, a parent or supervisor may confirm satisfactory completion of the tasks and then actuate the master switch to close the master circuit, thereby allowing transfer of transmissions from the transmission source to the electronic device such that the electronic device is operational.\nOther features, objects and advantages of the present invention will become apparent from a reading of the following description as well as a study of the appended drawings."}
{"text": "The technological advancement in the last two decades, have allowed for an evolution in the acquisition of goods and/or services, both physical and digital which has facilitated its purchase and delivery, regardless of the geographical location of the buyer user. For example, currently, the acquisition of various electronic goods such as music, videos, films, series and/or television programs, etc., can be done through the purchase of prepaid cards, which are acquired through a point of sale terminal of some establishment, where said prepaid card has a code that allows it the user to access the electronic good in the portal of the corresponding company.\nAn example of the above can be found in European Patent Application No. EP 1 519 332 A1, which discloses a method and system for providing digital content to consumers using a prepaid digital content medium. The method includes receiving unique identification information associated with the prepaid digital content medium (PDCM) from a PDCM distributor. Once a buyer has acquired a PDCM, the unique identification information associated with said PDCM is sent to an authentication service provider and the PDCM is activated based on said identification information. The authentication service provider sends an activation confirmation of the PDCM to the digital content provider in response to a confirmation request from said digital content provider, wherein the digital content provider sends a confirmation request in response to a request of digital content by a consumer who has acquired the PDCM.\nThe system above described gives the buyer the possibility to dispose and use the electronic good immediately after the transaction with the digital content provider is completed. In one of its modalities, said system uses prepaid cards that the user acquires in a physical establishment, which constitutes a limitation for the sales of the establishment, since the establishment has a “limited” number of cards to sell and in the case in which they are terminated, the user could not acquire the electronic good in said establishment. Therefore, the establishment must keep in physical stock, the prepaid cards of the digital content provider for sale to the public.\nIn the state of the art, we can find, in addition to the above prepaid cards, the use of printed codes in a purchase ticket, (known as “pin on receipt”), which is another way in which a user can purchase a code to access digital content from a provider. An example of this is described in U.S. Pat. No. 7,613,237, which comprises a system and method for providing a personal identification number (PIN) printed on a receipt by a customer terminal. The method includes storing at the client terminal, receipt formats that are sent from a server for various prepaid services from a plurality of prepaid service providers. The method further includes receiving, at the client terminal, a request for a PIN associated with one or more of the various prepaid services. The requested PIN and other information associated with one or more of the various prepaid services is displayed or printed according to one or more of the receipt formats for its reception and dispensing by the customer terminal.\nThe method above described eliminates the use of prepaid cards, which entails that the trade in question is no longer limited to the physical stock of said cards and has “unlimited” availability for the sale of codes to the user.\nHowever, said method is limited to the user moving to the location where the point of sale terminal is located. In addition, there is a risk that the user may damage or misplace the receipt where the PIN is printed, and in case this happens, it is not possible for the user to receive the PIN printed again at the establishment where he acquired it. Another disadvantage of the above described method is that there is a need for a printer, whether thermal, dot matrix, inkjet or laser, whereby the ability to generate sales in the establishments is limited to having a device that allows printing the access code and instructions for use.\nThe technological advance and the decrease in the costs of physical electronic products such as cell phones and tablets, have allowed its massification in our society, which means that most of the people currently have an equipment and assigned cell phone number. That is why many companies and organizations since a few years began to develop digital systems that allow the sending of notifications to the user cell phone number, which may contain information on an action or a code.\nParticularly, some banks have the service of sending the personal identification number (PIN) to their users via SMS type electronic messages, which prevents the user from having to physically move to the branch of the bank to obtain the PIN. With this, the user can promptly have his PIN for activation of a credit or debit card. However, this system is limited to sending these numbers for specific use in the services and products offered by the issuing bank.\nAccordingly, in the state of the art, to date there is no efficient system that allows the acquisition of electronic products and/or services by means of a code that is sent to a user device through an electronic message, which allows the user to have immediate access to the product and/or service, when acquired through the server portal of the provider."}
{"text": "The present invention relates to a double acting hydraulic piston and cylinder assembly and a method for assembling the same.\nDouble acting, double ended cylinders include a piston rod which extends from both ends of the cylinder. The piston operates in both directions by introducing fluid under pressure alternatively to opposite sides of the piston.\nIn the past, pistons for such double ended cylinders have been fastened to the rod by several different methods. One method comprises threading a two piece rod together within the piston so that the rod extends from opposite ends of the piston. Another method involves welding the piston in place on the rod. A third method involves installing collars on both sides of the piston for holding the piston against axial movement on the rod. A fourth method involves pinning the piston to the rod. Therefore, a primary object of the present invention is the provision of an improved double acting hydraulic piston and cylinder assembly and an improved method for assembling same.\nA further object of the present invention is the provision of means for securing the piston to the rod which eliminates the disadvantage of distortion of the rod which occurs in the present method of welding the piston to the rod.\nA further object of the present invention is the provision of a device for attaching the piston to the rod so that the piston can be fitted loosely over the rod and then attached securely in place.\nA further object of the present invention is the provision of a means for holding the piston positively in place once it is installed so that it will not move in either direction axially on the rod.\nA further object of the present invention is the provision of a device which may be disassembled easily for removing the piston and for repacking or rebuilding the cylinder and piston assembly.\nA further object of the present invention is the provision of an improved hydraulic piston and cylinder assembly wherein the piston rod is straight and completely concentric since there is only one part to the piston rod rather than the prior method of threading two piston rod segments into the piston.\nA further object of the present invention is the provision of a hydraulic piston and cylinder assembly which has interchangeable parts adapted to high production techniques."}
{"text": "During the implantation of an endocardial lead body, the lead is introduced into the heart using a venous approach, usually from the subclavian or cephalic vein in the shoulder area under the pectoral muscle. To stabilize the lead body at the venous entry site, the lead body is secured to both the vein and to the surrounding fascia tissue. A suture placed around the vein near the lead entry point ties the lead body to the vein, and a suture sleeve around the lead body is used to anchor the lead body to adjacent tissue.\nSuture sleeves in present use are generally tubular structures molded out of a soft, implantable elastomer such as silicone. After the lead body is tied to the vein, the sleeve is slid along the lead body to the location at which the lead is to be anchored to the underlying tissue. One or more sutures are then tied around the sleeve to compress it and thereby secure it to the lead body. Circumferential grooves in the outer surface of the sleeve are typically provided for this purpose. The last step is to anchor the sleeve to adjacent body tissue; sutures passed through eyelets formed in a pair of tabs projecting from the sleeve provide the required anchoring.\nThese existing suture sleeves have several drawbacks. For example, it is difficult for the physician to control the degree to which these sleeves are compressed when they are secured to the lead body. The ligature around the sleeve must be tight enough to prevent the lead body from sliding in the suture sleeve but not so tight as to damage the insulation of the lead body. This is especially important with bipolar coaxial leads because an excessively tight ligature could rupture the lead insulation and cause the outer and inner electrical leads to come into contact with each other, resulting in a short circuit. Overtightened ligatures can also result in electrical lead fractures. These problems are common enough to warrant the inclusion of cautionary information in pacemaker product literature or in notices included in the product packaging regarding the use of anchoring sleeves. It would therefore be desirable to eliminate the need for sutures for compressing the sleeve.\nU.S. Pat. No. 4,672,979 discloses another type of suture sleeve comprising an outer, elastomeric tubular sleeve and a relatively stiff collet member. The collet includes a plurality of axially extending legs having inner surfaces that are serrated so as to firmly grip the lead body. In use, the collet and the tubular sleeve are separately threaded onto the lead body. With the collet positioned at the lead body anchoring site, the tubular sleeve is forced over the collet thereby urging the serrated surfaces of the collet inwardly into engagement with the lead body. The assembly is then sutured to surrounding tissue to anchor the lead body.\nThe suture sleeve of the '979 patent has several disadvantages. The serrations have a tendency to pinch and puncture the outer insulative layer of the lead body. Moreover, because the tubular sleeve and collet are separate elements, during positioning of the collet the tubular sleeve may slide down the lead body due to its lubricity when covered with body fluids. In this event, the suture sleeve installation is substantially complicated.\nAccordingly, it is an object of the present invention to provide a suture sleeve for securely gripping and anchoring the lead body of an implantable medical device, such as a cardiac pacemaker, that does not require the use of sutures to compress the sleeve.\nIt is another object of the present invention to provide a lead body gripping and anchoring suture sleeve that does not damage the lead body.\nIt is yet another object of the present invention to provide a compressible lead body gripping and anchoring suture sleeve in which the compression of the sleeve is self-limiting so as to prevent lead body damage."}
{"text": "Many satellites and other systems have one or more covers or doors to protect internal mechanisms, e.g., optics, sensors and sensitive machinery during transport or other periods of disuse. These types of doors utilize various mechanisms, such as mechanical launch locks to hold them in a closed position and motors to open them. Motors and the control circuitry by which they are operated are expensive to design, procure, manufacture and test. In some applications, such as space-based systems, testing typically includes a life cycle test which could take months to complete. Additional items that drive cost in such space-based systems include the launch lock, the control/telemetry electronics, the drive motor, the gear box, as well as vacuum compatible bearings and grease. Research and development efforts that seek to simplify cover/door mechanisms, and to reduce cost accordingly, are ongoing."}
{"text": "Thermoelectric heat transfer systems are well known for use in refrigeration. A typical thermoelectric heat transfer system includes at least one thermoelectric module to create a temperature differential. More specifically, when energized, heat moves across the thermoelectric module to form a hot surface and a cold surface. The cold surface provides the cooling needed for refrigeration.\nIn recent years, improvements have been made to utilize coolant circuits to draw heat off of the hot surface of the thermoelectric module to further improve the efficiency of refrigeration. Referring to U.S. Pat. No. 5,653,111 to Attey et al., a thermoelectric heat transfer system utilizing a coolant circuit for this purpose is shown. In Attey et al., the thermoelectric heat transfer system includes a pump in fluid communication with the coolant circuit to circulate coolant. A first heat exchanger is disposed in fluid communication with the coolant circuit to remove heat from the coolant being circulated. At the same time, a manifold is in fluid communication with the coolant circuit. An outer surface of the manifold is in contact with a hot surface of a thermoelectric module. By thermally connecting the coolant circuit with the hot surface of the thermoelectric module, the coolant can draw heat from the hot surface to improve the cooling efficiency of the thermoelectric module. The cold surface of the thermoelectric module is in contact with an outer surface of a second manifold to cool fluid flowing through the second manifold.\nThermoelectric heat transfer systems for cooling beverage containers are also well known in the art. A typical heat transfer system for cooling a beverage container includes a sleeve adapted to receive the beverage container. In these systems, the thermoelectric module is disposed within the sleeve to create a temperature differential between the sleeve and the beverage container. More specifically, a hot surface of the thermoelectric module is in contact with the sleeve, while a cold surface of the thermoelectric module is in contact with the beverage container thereby drawing heat from the beverage container. A fan assembly draws heat away from the sleeve. An example of such a system is shown in U.S. Pat. No. 6,530,232 to Kitchens."}
{"text": "1. Field of the Invention\nThe present invention relates to a selective catalytic process for preparing N-halothiosulfonamide modified terpolymers of ethylene-alpha olefin-diene, and catalyst compositions suitable for use in this process.\n2. Description of Information Disclosures\nModified polymers of the type of interest herein, specifically rubbery polymers, are described in U.S. Pat. No. 3,915,907 and U.S. Pat. No. 3,970,133 (R. J. Hopper) incorporated herein by reference. The desirability and advantages of such polymers was additionally disclosed in articles authored and co-authored by R. J. Hopper (Rubber Chemistry and Technology, Vol. 49, p. 341-352 (1976) and International Conference on Advances in the Stabilization and Controlled Degradation of Polymers, Lucerne, Switzerland, May 23-25, 1984).\nThese references describe the reaction of various N-chlorothioamides, and particularly N-chlorothio-sufonamides with various ethylene propylene diene monomer (EPDM) rubbers in solutions (see, e.g., Examples I-XII of U.S. Pat. No. 3,915,907) as well as by mixing in a Banbury.RTM. type internal mixer (Examples XIII-XVI of U.S. Pat. No. 3,915,907). Furthermore, in Example XVII of the same reference, there is disclosed in a solution reaction the use of powdered anhydrous zinc chloride (this is referred to as a catalyst in Table I, footnote g and page 346 of the cited Rubber Chemistry and Technology article). The references contain no further disclosure concerning the effect of zinc chloride nor do they suggest the use of other catalytic agents.\nU.S. Pat. No. 4,820,780 issued Apr. 11, 1989 discloses a method of preparing a N-chlorothiosulfonamide-modified terpolymer of ethylene, alpha olefin containing from 3 to 6 carbon atoms and a diene in the presence of a saturated aliphatic carboxylic acid containing from 6 to 30 carbon atoms, such as stearic acid. Calcium carbonate may also be added to the reaction as a scavenger for the HCl by-product. It does not suggest the use of metal salts of weak acids as catalysts.\nIt would be desirable to improve the rate of reaction, the conversion and product quality produced by known processes to achieve commercially useful conditions. It would be particularly useful to improve the melt phase process so as to enhance the addition reaction while minimizing the crosslinking reaction.\nIt has now been found that certain catalyst compositions comprising zinc salts or iron salts of weak acid in combination with a metal salt of certain metals (other than zinc or iron) of weak acids (i.e., co-catalyst) are selective catalysts in the preparation of N-halothiosulfonamide modified terpolymers of ethylene-alpha olefin-diene. By \"selectively,\" it is intended herein that the desired addition reaction is increased while the undesired crosslinking reaction, which leads to increased Mooney viscosity and to gelation, is minimized.\nMixtures of blends of low-unsaturation rubbery polymers, particularly EPDM, with highly unsaturated rubbery polymers are of practical importance because of the superior ozone resistance imparted to the blend by the low-unsaturation EPDM rubber. Unfortunately, the presence of the low-unsaturation rubber also affects the mechanical and hysteresis characteristics of the vulcanizates in an adverse manner, as evidenced by lower tensile strength and modulus values, and by higher dynamic heat build-up and permanent set. These undesirable responses are contributed to by the mutual incompatibility of the two types of rubber, the substantially slower cure rate of the low-unsaturation rubber, and the greater affinity of the typically polar curvatures for the high-unsaturation rubber. As a result, the vulcanized blend is a heterogeneous dispersion of poorly cured and/or largely uncured low-unsaturation rubber in well cured high-unsaturation rubber. The resulting poorer mechanical and hysteresis properties severely limit, or preclude, the use of such blends in articles subjected to severe service requirements such as tires. A method for improving the physical and dynamic properties of such rubbery polymer vulcanized blends is of considerable commercial and practical importance."}
{"text": "Multiuser diversity is one of diversities caused by independently time-varying a channel environment between users in a wireless communication system. Multiuser diversity gain is obtained by allocating a resource when a channel of each user reaches the best condition through scheduling during observation of a change of a channel of each user. As the change of the channel is increased, the diversity gain is increased. Accordingly, the diversity gain is limited under a channel environment where the change of the channel is small or slow. Under the above-mentioned environment, the diversity gain is generated by inducing a large and quick change of the channel using multiple transmission antennas. This method is called opportunistic beamforming (hereinafter referred to as “OBF”) or random beamforming.\nOBF technology is limited to the diversity thereof under an environment corresponding to a small number of users, since it uses only one beam to search an optimal user at each time slot. Therefore, technologies that generate multibeams and select an optimal beam of each user to increase multiuser diversity gain, have been proposed in order to overcome the above-mentioned limitation. As these technologies in the related art, there are codebook-based opportunistic beamforming (hereinafter referred to as “COBF”) and codebook-based beamforming with code rotation (hereinafter referred to as “CBF-CR”), where the codebook contains multiple vectors used for multibeams.\nThe COBF uses a pilot only for channel estimation, and allows each user to estimate a channel. A base station receives from each user the maximum beam index and channel quality information (hereinafter referred to as “CQI”) corresponding to the beam, as feedback. Then, the base station performs scheduling by allocating a user having the maximum gain to each beam thereof, and transmits data. In the COBF method, a unitary matrix varying with time is multiplied by one of beam vectors in codebook, so that a beam varying with time is obtained. The COBF has a merit in that a pilot is fixed to a pilot used for channel estimation and a selective diversity effect is thus increased to infinity.\nThe CBF-CR generates multiple codebooks by grouping large-sized codebooks and then obtains multiple beams that vary with time by diversely using the codebooks with change of time. Therefore, the same effect as the COBF method is obtained in the CBF-CR.\nIn the COBF, multibeams have been obtained by multiplying unitary matrices, which vary with time, by beam vectors of one codebook. In contrast, in the CBF-CR, multibeams are obtained by generating multiple codebooks, and then using the codebook while exchanging the codebook with time.\nHowever, the method of designing a beam and the method of designing a beam parameter in the related art have a problem in that performance deteriorates.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "1. Field of the Invention\nThe invention relates to a decoding apparatus, more particularly to a Context-based Adaptive Binary Arithmetic Coding (CABAC) decoding apparatus and a decoding method thereof.\n2. Description of the Related Art\nH.264/AVC is a recently developed video signal compression standard, which employs advanced encoding techniques to provide superior compression efficiency compared to other existing video signal compression standards such as MPEG-1, MPEG-2, H.263, etc.\nH.264/AVC adopts Context-based Adaptive Variable Length Coding (CAVLC) technology or CABAC technology for entropy encoding to increase the coding efficiency. While the coding efficiency of CABAC technology is about 10% higher than that of CAVLC technology, encoding/decoding calculations in CABAC are more complex than those in CAVLC. During the CABAC decoding process, an input bit stream must at least go through an initialization step and an arithmetic decoding step before the decoding procedure can be completed. The total amount of time required by the initialization step and the arithmetic decoding step, which is actually relatively long, has a critical affect on the efficiency of a decoding system. As such, it is highly desired to reduce latency during the CABAC decoding process so as to improve the overall efficiency of the decoding system."}
{"text": "Circular extrusion dies with inner and outer lips are usually adjusted for concentricity by a circle of bolts located radially around the outer die lip. Concentricity adjustment is necessary to achieve a consistent transverse thickness of the material leaving the die. In most cases, the adjusting bolts move the complete outer lip to obtain concentricity with the inner lip. A compromise is needed due to temperature and machining variations around the two lips. Where a higher degree of control is desired, a flexible outer lip may be used. The adjusting bolts push on segmented shoes which flex a small sector of the outer lip. This forces the lip into the desired shape to achieve uniform thickness of material. Heated expansion bolts are a further sophistication allowing a finer, remote controlled adjustment. Flexible lip dies suffer from a number of deficiencies, such as undesired distortion of areas adjacent to the sector being flexed by an adjusting bolt, requirement for large forces to flex the lips making repeatability and fine adjustment difficult, necessity for construction of a very thin flex lip which is then susceptible to machining distortion and damage, and complications in heat transfer to and from the outer die lip due to thin construction and intimate contact with the segmented shoes around the periphery.\nHahn, et al., U.S. Pat. No. 4,548,570, describes an extrusion apparatus for producing thermoplastic pipe. Adjustment of the thickness of the pipe is achieved by varying the heat to heater cartridges which causes thermoplastic material in proximity to the heated section of sleeve to become less viscous and to flow at a faster rate. This increases the amount of thermoplastic material at that location, compared with lesser heated sections of the extrusion passageway. Additionally, radial repositioning of the sleeve due to expansion of the heater bolts will increase the width of the arcuate portion of the melt passageway in proximity to the heated section of the sleeve, causing increased thickness of the thermoplastic material. This heating operation does not appear to take place at the die lips. Phipps, U.S. Pat. No. 4,201,534 describes a foam extrusion die assembly which is adjusted by means of screws in the area of the die lips. Yazaki, et al., U.S. Pat. No. 4,057,385, describes cooling a deckle so as to increase the viscosity of the molten resin around the deckle, to prevent leakage. Solop, U.S. Pat. No. 3,984,508, also describes a method of controlling an extrudate using a deckle bar having heating and cooling means. Porter, et al., U.S. Pat. No. 3,973,890, describes adjusting relative positions of the die plates to maintain a constant pressure of heat plastified material. The temperature of the pellet cutting head can be heated or cooled by circulating fluids."}
{"text": "In this century, the shortage of fresh water will surpass the shortage of energy as a global concern for humanity; and these two challenges are inexorably linked, as explained, for example, in the “Special Report on Water” in the 20 May 2010 issue of The Economist. Fresh water is one of the most fundamental needs of humans and other organisms; each human needs to consume a minimum of about two liters per day. The world also faces greater freshwater demands from farming and industrial processes.\nThe hazards posed by insufficient water supplies are particularly acute. A shortage of fresh water may lead to a variety of crises, including famine, disease, death, forced mass migration, cross-region conflict/war, and collapsed ecosystems. Despite the criticality of the need for fresh water and the profound consequences of shortages, supplies of fresh water are particularly constrained. 97.5% of the water on Earth is salty, and about 70% of the remainder is locked up as ice (mostly in ice caps and glaciers), leaving only a fraction of all water on Earth as available fresh (non-saline) water.\nMoreover, the earth's water that is fresh and available is not evenly distributed. For example, heavily populated countries, such as India and China, have many regions that are subject to scarce supplies. Further still, the supply of fresh water is often seasonally inconsistent. Meanwhile, demands for fresh water are tightening across the globe. Reservoirs are drying up; aquifers are falling; rivers are dying; and glaciers and ice caps are retracting. Rising populations increase demand, as do shifts in farming and increased industrialization. Climate change poses even more threats in many regions. Consequently, the number of people facing water shortages is increasing. Naturally occurring fresh water, however, is typically confined to regional drainage basins; and transport of water is expensive and energy-intensive.\nAdditionally, water can be advantageously extracted from contaminated waste streams (e.g., from oil and gas production) both to produce fresh water and to concentrate and reduce the volume of the waste streams, thereby reducing pollution and contamination and reducing costs.\nNevertheless, many of the existing processes for producing fresh water from seawater (or from brackish water or contaminated waste streams) require massive amounts of energy. Reverse osmosis (RO) is currently the leading desalination technology. In large-scale plants, the specific electricity required can be as low as 4 kWh/m3 at 30% recovery, compared to the theoretical minimum of around 1 kWh/m3; smaller-scale RO systems (e.g., aboard ships) are less efficient.\nOther existing seawater desalination systems include thermal-energy-based multi-stage flash (MSF) distillation, and multi-effect distillation (MED), both of which are energy- and capital-intensive processes. In MSF and MED systems, however, the maximum brine temperature and the maximum temperature of the heat input are limited in order to avoid calcium sulfate, magnesium hydroxide and calcium carbonate precipitation, which leads to the formation of soft and hard scale on the heat transfer equipment.\nHumidification-dehumidification (HDH) desalination systems include a humidifier and a condenser as their main components and use a carrier gas (e.g., air) to desalinate brine streams. A simple version of this technology includes a humidifier, a condenser, and a heater to heat the brine stream. In the humidifier, hot brine comes in direct contact with dry air, and this air becomes heated and humidified. In the condenser, the heated and humidified air is brought into (indirect) contact with a coolant (for example, cold brine) and gets dehumidified, producing pure water and dehumidified air. The HDH process operates at lower top brine temperatures than MSF and MED systems, precipitation of scaling components is hence avoided to some extent.\nAnother approach, described in U.S. Pat. No. 8,119,007 B2 (A. Bajpayee, et al.), uses directional solvent that directionally dissolves water but does not dissolve salt. The directional solvent is heated to dissolve water from a salt solution into the directional solvent. The remaining highly concentrated salt water is removed, and the solution of directional solvent and water is cooled to precipitate substantially pure water out of the solution.\nThe present inventor was also named as one of the inventors on the following patent applications that include additional discussion of HDH and other processes for purifying water: U.S. application Ser. No. 12/554,726, filed 4 Sep. 2009; U.S. application Ser. No. 12/573,221, filed 5 Oct. 2009; U.S. application Ser. No. 13/028,170, filed 15 Feb. 2011; and U.S. application Ser. No. 13/241,907, filed 23 Sep. 2011; U.S. application Ser. No. 13/550,094, filed 16 Jul. 2012; U.S. application Ser. No. 13/916,038, filed 12 Jun. 2013; and U.S. application Ser. No. 13/958,968, filed 5 Aug. 2013."}
{"text": "The present invention relates to a process for manufacturing bodies of polymeric material by means of the blow-molding technique, in particular to provide bodies comprising a core of foamed material enclosed by an external shell.\nSuch bodies normally display a rigid structure and good properties of heat insulation, so to make the use thereof for such household electrical appliances, as refrigerators, and the like, particularly advantageous. At present, the technique of blow-molding is used in order to manufacture a hollow shell of thermoplastic material, inside which a thermosetting foam is subsequently injected, for example a polyurethane foam obtained by means of the well-known RIM technology. Such an operating way shows problems both in terms of possibility of finished article recycling, and in terms of environmental pollution by the foaming agent.\nIn fact, the chemical nature of the foamed polyurethanic core makes it impossible for the manufactured article to be recycled, and the chloro-fluoro-alkane foaming agent normally used in RIM technology causes problems of interaction with ozone in the stratosphere. Furthermore, the execution of the process according to the two-step technology known from the prior art (blow-molding and subsequent foam injection) causes the production cycle to be rather slow and expensive.\nThe purpose of the present invention is of supplying a process of the type as specified in the introduction to the specification, which does not display these drawbacks and which is simple and inexpensive to carry out.\nAccording to the invention, such a purpose is achieved because the blowing is performed by introducing into the interior of a hollow thermoplastic body a foamable thermoplastic polymeric material suitable for bringing the walls of the hollow body into contact with the walls of a mold.\nBy \"foamable thermoplastic polymeric material\" any thermoplastic material mixed with a foaming agent, e.g., CO.sub.2, N.sub.2, hydrocarbons, are meant. It is evident that, as in the traditional blow-molding technology, the hollow extruded body (parison) is at a temperature close to the polymer softening temperature, and anyway such as to secure the \"blowing\" of the hollow body. This feature is particularly advantageous in the case of the present invention, because the foamed thermoplastic material (which also is at a temperature which is close to the polymer softening temperature) can be thermowelded to the walls of the hollow body, thus making it possible to obtain a structurally resistant manufactured article. A further advantage of the present invention relates to the possibility of using the same polymeric material, or materials belonging to the same polymeric family, for both the foamed core and the external shell, making it possible for the manufactured article to be recycled at the end of its useful life time. The process according to the present invention also makes it possible for the production cycle to be speeded up, as compared to the traditional technology of blow-molding and subsequent foaming of the cooled hollow shell.\nA further purpose of the present invention is providing a device for carrying out the above said process.\nSuch a device, of the type which comprises an extruder with an accumulation chamber suitable for cyclically extruding a tubular parison between two mobile half-molds, is characterized in that said device also comprises a second extruder for a polymeric material admixed with a foaming agent, and an auxiliary accumulation chamber in communication with said second extruder, and equipped with an injection nozzle centrally installed relative to a ring-shaped nozzle for the tubular parison extrusion.\nFurther advantages and characteristics of the process and device according to the invention will be clear from the following disclosure in detail, supplied for merely exemplifying, non-limitative, purposes, by referring to the accompanying drawings, in which:"}
{"text": "Air intake passages for automobile heating and air conditioning systems have been known since the advent of climate control systems for passenger compartments of automobiles. With the growing popularity of vans and other vehicles with larger-than-standard passenger compartments, there is an increased need to provide climate control systems which will efficiently control the climate in the larger passenger compartment.\nWhile it has been known that a larger air flow through a climate control system will increase the efficiency and rate at which a system can process, and thus control, the volume of air in the passenger compartment, attempts to increase the air flow into climate control systems have had limited success. The application of the principle of achieving a more efficient heat exchange by increasing the air flow to practical use for a motor vehicle passenger climate control system has not been an easy matter. In large part, this is because the spatial constraints of the application do not lend themselves to simply increasing the cross-section of the intake passage.\nOther prior art approaches to increasing air flow have included providing additional entries into the intake passage of the vehicle climate control system. For example, Japanese Patent 58-101813 to Yamashita provides a second door 17 to an inlet 18, which door is interlocked to a door 2 of another intake 1. This approach merely provides two doors in the same plane of the housing wall, requiring a larger cross-sectional area to be available in that plane.\nIn contrast, U.S. Pat. No. 4,802,405 discloses three free-circulated air openings in three different walls of its system casing, but each inlet requires control by a separate damper door. U.S. Pat. No. 4,779,672 similarly discloses a system having two recirculation air inlets, each of which is provided with its own damper door, with corresponding control mechanisms.\nFrench Patent 2,556,289 discloses a mixer box for a vehicle interior heater which provides for selective mix of partial recirculation air which comes entirely from window openings provided in a single wall.\nU.S. Pat. No. 4,676,272 generally exemplifies an air flow control passage also not dedicated to the problem of increasing air flow of recirculated air into a vehicle climate control system. Similarly, the system of German Patent 3,421,323 discloses an air distributor which does not address increasing air flow of recirculation air into a system by multiple entries in different walls of the casing.\nAs will be seen below, the multiple entry air intake passage of the present invention provides a simplified design for increasing air flow while simultaneously eliminating the need for extra doors and door actuating mechanisms, in turn providing savings in cost and size."}
{"text": "1. Field Of The Invention\nThis invention relates to computer systems and, more particularly, to methods and apparatus for automatically connecting computers to one of two available local area networks (LANs).\n2. History Of The Prior Art\nOne of the significant movements in the computer industry has been toward interconnectivity of individual computers by local area networks. Initially, the desire was to allow the sharing of underused assets such as printers. More recently, the impetus has been the desire to share application programs and data and to enhance communication.\nWhatever the reasons for their existence, LANs are proliferating. Typically, a LAN will be installed initially to provide connection between a small group of local computer users. The first LANs which were installed utilized coaxial cable as the physical wiring for the transmission of signals. Coaxial cable was necessary in order to reduce the noise and allow operation in the face of signal attenuation. As technology advanced and standards changed, it became possible to utilize twisted pair wiring as the physical connection to join the computers of a LAN together. Twisted pair wiring is preferred because it is less expensive and less expensive to install.\nOften it is desirable to change computers from a coaxial network to a twisted pair network. During some interim period, it may be desirable to connect the computers to both networks. In a similar manner, it may become desirable to connect computers already joined in a local twisted pair network to another network in which the computers are joined by coaxial cable. In either case, it is possible that a number of individual computers will be physically connected simultaneously to two separate networks, a twisted pair network and a coaxial cable network. It is desirable that these computers be able to utilize either one of the two networks which is available and functioning so that if one network is not functioning communications will automatically be established on the other network."}
{"text": "Digital content distribution is extremely popular. Most content producers desire some sort of rights management, and so content distribution systems typically include content management systems that control aspects of playback, copying, or other usage of the digital content. As one example, a device in the content distribution chain may be configured to enforce terms of a license provided by a content producer or other authority. The license is stored as one or more computer-accessible files or other data structures accessed by the device to enable (or deny) access or other use of the content."}
{"text": "As consumer devices have gotten smaller and smaller in response to consumer demand, the individual components of these devices have necessarily decreased in size as well. Semiconductor devices, which make up a major component of devices such as mobile phones, computer tablets, and the like, have been pressured to become smaller and smaller, with a corresponding pressure on the individual devices (e.g., transistors, resistors, capacitors, etc.) within the semiconductor devices to also be reduced in size.\nOne enabling technology that is used in the manufacturing processes of semiconductor devices is the use of photolithographic materials. Such materials are applied to a surface and then exposed to an energy that has itself been patterned. Such an exposure modifies the chemical and physical properties of the exposed regions of the photolithographic material. This modification, along with the lack of modification in regions of the photolithographic material that were not exposed, can be exploited to remove one region without removing the other.\nHowever, as the size of individual devices has decreased, process windows for photolithographic processing as become tighter and tighter. As such, advances in the field of photolithographic processing have been necessitated in order to keep up the ability to scale down the devices, and further improvements are needed in order to meet the desired design criteria such that the march towards smaller and smaller components may be maintained.\nCorresponding numerals and symbols in the different figures generally refer to corresponding parts unless otherwise indicated. The figures are drawn to clearly illustrate the relevant aspects of the embodiments and are not necessarily drawn to scale."}
{"text": "The present invention relates to the field of integrated circuit devices. More specifically, it relates to testing and sorting individual integrated circuit devices according to their operational characteristics.\nIntegrated circuit devices, or dies, are manufactured through a multi-step process during which transistors and other electronic circuit components are formed in areas on a thin silicon wafer. After the circuit fabrication is completed, the individual dies are cut from the silicon wafer.\nThe manufacturing process for integrated circuits is imprecise; minor variations during each step of the manufacturing process affect vital performance characteristics of individual dies. These variations impact the speed at which a particular die operates, which directly impacts the commercial value of each individual die. It is therefore desirable to the manufacturer to determine the speed at which individual dies can operate so the manufacturer can then sort them accordingly.\nTypical integrated circuit test devices consist of bench setups utilizing standard lab equipment including an oscilloscope or frequency meter to directly or indirectly measure the operating frequency (speed) of the circuit. In some setups such outboard test equipment is used to test the speed of a device, such as an oscillator, which is fabricated in association with a die but which is not materially part of the operational die circuit. The tested frequency of the oscillator is representative of the operational speed of the die. All these systems, however, require expensive specialized and sometimes bulky frequency measurement equipment to be provided with the test setup.\nWhat is needed is a simplified method and apparatus for determining the operating speed of the die.\nThe present invention provides a method and apparatus for characterizing an integrated circuit device based on device switching speed on an individual die or lot-by-lot basis by measuring frequency in a digital domain. An oscillator, integrated in association with the die and capable of generating pulses, and a counter are used during a test period to count the pulses from the oscillator to a predetermined value. A pulse generator is initiated at the beginning of a test which resets the counter to a known state and starts the oscillator. When the counter reaches its predetermined maximum count value a signal is issued. The elapsed time between the initiation of the oscillator and when the counter full signal is issued represents the summation of many oscillator periods. The period of the oscillator may then be determined by the formula: oscillator period=elapsed time/capacity of counter. The determined oscillator period provides an indication of the operating speed of electronic devices on the die. A digital tester which performs other tests on the die can be used to record the elapsed time and thereby determine the operating speed of a die.\nThese and other features of the invention will be more clearly understood from the following detailed description of the invention which is provided in connection with the accompanying drawing."}
{"text": "The present invention relates to firearms.\nIn particular, the present invention relates to that type of firearm, such as a semi-automatic shotgun, which has a bolt capable of moving axially between a forward firing position and a rear retracted position as well as capable of rotating between locked and unlocked positions. Such a bolt, which forms part of the action of the firearm, is housed within a receiver which is formed with an ejection port through which a cartridge case of a spent cartridge is ejected. The movement of the bolt is brought about at least partially by way of a sleeve which surrounds the bolt, and in addition to its other functions, this sleeve serves to close the ejection port when the bolt is in its forward firing position.\nBecause the ejection port is relatively long and because it must be covered by this sleeve, the sleeve also is undesirably long, and in addition the receiver must be made of a length sufficiently great to receive the relatively long sleeve when the bolt is retracted in order to uncover the ejection port at least to an extent sufficient to permit the cartridge case to be ejected.\nThus, these factors result in an undesirably long action and receiver."}
{"text": "The present invention relates to a method for analyzing a stereoscopic image of a fundus to detect and analyze the optic disk on the basis of the stereoscopic image of the fundus, and an apparatus for executing that method.\nGlaucoma is known as a disease of the visual function. In the examination of glaucoma, it is said that the observation of patient's optic disk on the basis of an image of the fundus is important. In recent years, it has become possible to quantitatively measure the optic disk due to progress in the technology of image analysis. As the results of analysis, the following parameters are computed: the C/D ratio (a distance ratio between a rim of an optic disk cupping (the optic disk cupping being referred to as a cup) and a rim of the optic disk (the optic disk being referred to as a disk)), the area of the plane within the cup rim, the depth of the cup, the volume of the cup, and the like.\nThe rim of the cup which serves as a basis for the computation of these parameters is conventionally detected such that, by using the disk as a reference, positions which are lowered by an arbitrary depth (e.g., 150 .mu.m) from the disk are estimated to plot the rim of the cup.\nHowever, since it is not uncommon that overall disk regions of the patent's eyes are inclined, there arise cases where the rim of the cup cannot be detected by the above-described detection method.\nIt is possible to obviate the above-described problem if an operator designates three points on the disk region and effects analysis by using as a reference plane an inclined plane thus determined. In accordance with this method, however, there is a problem in that the data of the cups detected vary for each operator, and the data can be different on each occasion of analysis even for the same operator, with the result that the reliability and reproducibility of the results of analysis are poor."}
{"text": "Non-volatile semiconductor memories (NVSMs) are widely used in many electronic devices such as personal digital assistants (PDAs), laptop computers, mobile phones and digital cameras. Some of these memories have arrays of charge storage transistors, such as floating gate transistors."}
{"text": "This invention relates to a blood collection assembly incorporating a microcollection container, and is related to the subject matter disclosed in co-pending application Ser. No. 695,121, filed Jan. 28, 1985. The invention is an improvement over the collection assembly described and claimed in U.S. Pat. No. 4,397,318, issued Aug. 9, 1983, which is hereby incorporated by reference in its entirety. Reference should be made to that patent for background information concerning the teachings of the invention here. The earlier patent involved the use of a scoop collector for connection to a blood microcollection container for engaging a puncture wound to obtain a blood sample from an individual for subsequent examination of that sample for the determination of the presence or absence of some disease or other problem in a patient. The scoop-type blood collection device provides a substantially larger engaging surface for engaging the puncture for collecting the blood, and a substantially larger transfer surface for rapidly transferring the blood from the collector into the microcollection container. Because of the relatively large engaging surface for engaging the puncture wound, the arrangement does not require a precise positioning of the scoop engaging surface in order to initiate and rapidly transfer a quantity of blood to the microcollection container.\nAs will be appreciated by practitioners-in-the-art, recent advancements in analytical instrumentation have made it possible to carry out a variety of hematological or chemical diagnostic procedures on very small quantities of blood. Because of this, a patient's heel, finger or earlobe may be punctured and a very small quantity of blood rapidly collected into a microcollection container for such testing. Such arrangements obviate the need to withdraw venous blood from patients. However, such collection arrangements must be such that the blood is rapidly collected prior to any coagulation thereof. In the past, prior to the scoop collector disclosed in the above-noted U.S. Pat. No. 4,397,318, a cap or top arrangement was configured to fit on the top of a microcollection container with the top having an integral capillary tube for engaging the puncture and transferring blood to the container. However, with such an arrangement, the tip of the capillary tube had to be arranged precisely adjacent the puncture wound and the entire apparatus had to be so positioned that the blood flow along the bottom surface of the tubular microcollection container moved continuously in order to engage the surface of the container. Otherwise, if a precise positioning was not carried out, capillary action was not initiated or was slowed with subsequent clotting. Representative such collectors are taught in U.S. Pat. No. 4,024,857, issued May 24, 1977.\nOne problem with the scoop collector taught and claimed in U.S. Pat. No. 4,397,318, although the arrangement taught therein is highly efficient for the rapid collection of a blood sample into a microcollection container, is the fact that the assembly for making the collection must be distributed with the microcollection container or tube having a separate cap. Typically, the technician must remove the cap, and place on the container the scoop collector prior to making a collection of a blood sample. Subsequent to this collection, moreover, the scoop collector must then be removed, and the cap replaced on the container for delivery to a lab for investigation of the sample. Such removal and replacement of parts on the top of the blood microcollection container is cumbersome, as one will understand, particularly if the technician is, for example, attempting to take a blood sample from a screaming, wiggly baby. Moreover, the technician or nurse or doctor may become exposed to the blood sample during this transfer procedure in removing the blood collection scoop arrangeent and replacing the cap on the blood microcollection container.\nWith the invention claimed in the above-noted co-pending application Ser. No. 695,121, by contrast, a scoop arrangement is incorporated into a blood microcollection assembly in such a way that the scoop collector does not have to be removed until such time as the technician in the laboratory wishes to obtain access to the sample in the blood microcollection container. This is achieved by the use of a cap which is a two-position cap. That is, the assembly is distributed to potential users with the cap in place over the scoop collector on the top of the blood microcollection container.\nWhen the nurse or doctor wishes to take a blood sample, the cap is removed and the front end of the scoop collector is placed adacent the wound for collection of blood. Once the blood sample has been taken, the cap is again placed over the scoop collector without any removal of the scoop collector, as in the past. Then, the technician merely has to press-fit the cap down over the scoop collector. This press-fit movement has the effect of permanently locking the cap onto the scoop collector. Therefore, access to the blood sample in the container cannot be obtained unless the cap and the scoop collector arrangement are removed simultaneously. For this reason, no one can be exposed to any blood left in or around the scoop collector arrangement after the sample has been taken, and until such time as the sample is to be obtained from the microcollection tube at the lab.\nAs practitioners-in-the-art will understand, this arrangement reduces the amount of fumbling and movements necessary during the course of taking a blood sample while at the same time reducing the possibility of contamination to the nurse or anyone else present during the taking of the sample.\nWith the invention claimed in this application, any difficulty engendered in taking a blood sample is further reduced dramatically by providing in the assembly itself a built-in lancet for making the wound for collecting the blood sample. That is, the lancet is positioned in the blood collector of the assembly and covered by the cap of the assembly prior to use. This allows for sterilization of the lancet and preservation of the sterile condition prior to use.\nWhen a sample of blood is to be collected with the combination assembly of the invention, the cap is removed. Then, the nurse or technician has an exposed collector and an exposed lancet simultaneously. The wound is made with the lancet and the sample collection made without the laying down or picking up of separate devices. The technician's hand is in place at the site of the wound for immediate collection.\nOnce collection of the sample is completed, the cap is again placed over the collector of the assembly, and moved or forced into its second locking position, as described in the above-noted co-pending application. By this capping procedure, not only is the blood collector removed from contaminating anyone coming into contact, but also, in accordance with this invention, the lancet itself is also capped and removed from such exposure. They are, as will be appreciated by a review of the co-pending application, permanently locked under the cap placed over the blood collection assembly.\nOther advantages of the invention include the fact that but a single package must be made up and sterilized to achieve collection of a blood sample. Moreover, only this one package must be opened and handled to obtain collection of the sample followed by proper disposal of contaminating objects.\nOther objects and advantages of this invention will be apparent from the following description, the accompanying drawings and the appended claims."}
{"text": "Small diameter glass fibers are useful in a variety of applications including acoustical or thermal insulation materials. When these small diameter glass fibers are properly assembled into a lattice or web, commonly called a wool pack, glass fibers which individually lack strength or stiffness can be formed into a product which is quite strong. The glass fiber insulation which is produced is lightweight, highly compressible and resilient. For purposes of this patent specification, in using the terms \"glass fibers\" and \"glass compositions\", \"glass\" is intended to include any of the glassy forms of mineral materials, such as rock, slag and basalt, as well as traditional glasses.\nThe common prior art methods for producing glass fiber insulation products involve producing glass fibers from a rotary process. A single molten glass composition is forced through the orifices in the outer wall of a centrifuge or spinner, producing primarily straight glass fibers. The fibers are drawn downward by a blower, and conventional air knife and lapping techniques are typically used to disperse the veil. The binder required to bond the fibers into a wool product is sprayed onto the fibers as they are drawn downward. The fibers are then collected and formed into a wool pack. The wool pack is further processed into insulation products by heating in an oven, and mechanically shaping and cutting the wool pack.\nIdeally, insulation products of glass fibers would have uniform spacing between fibers assembled in the lattice. Glass fiber insulation is basically a lattice which traps air between the fibers and prevents circulation of air to inhibit heat transfer. As well, the lattice also retards heat transfer by scattering thermal radiation. A more uniform spacing of fibers would maximize scattering and, therefore, have greater insulating capability.\nIn the production of wool insulating materials of glass fibers, it becomes necessary to use fibers that are relatively short to achieve desirable lattice properties. Known lapping techniques for dispersion of short fibers in a veil have provided acceptable, although not ideal fiber distribution. By contrast, long fibers tend to become entangled with each other, forming ropes or strings. For purposes of this patent specification, in using the terms \"short fibers\" and \"long fibers\", the term \"short fibers\" is intended to include fibers of approximately 2.54 centimeters (approximately 1 inch) and less, and \"long fibers\" are intended to include fibers longer than approximately 5.08 centimeters (approximately 2 inches).\nLong fibers are more prone to entangle than short fibers, due, in part to their different aerodynamic properties, in addition to fiberizer throughput and geometry. Moreover, the longer they are, the more the long fibers tend to entangle. Conventional lapping techniques have failed to eliminate, and rather tend to enhance, formation of ropes and strings in veils of long or semi-continuous fibers. Even when undisturbed, veils of long fibers tend to form ropes and strings as the veil slows in its descent to the collection surface. Despite movement of the collection surface, long glass fibers (as do undisturbed veils of short fibers) tend to pile up into nonuniform packs of fibers, and unmanageable fiber accumulations. These nonuniform packs, characterized in part by roping and string formation, have long prevented significant commercial use of long fibers. The ropes of long fibers produce a commercially undesirable appearance and, more importantly, create deviation from the ideal uniform lattice and reduce the insulating abilities of the glass wool.\nHowever, even short fibers that are straight form only a haphazard lattice, and some of the fibers lie bunched together. As a result, existing glass wool insulating materials continue to have significant non-uniformities in the distribution of fibers within the product. Thus, the ideal uniform lattice structure cannot be achieved.\nA further problem presented by use of short straight fibers is the binder material necessarily added to the fibers to provide product integrity. Binder provides bonding at the fiber to fiber intersections in the lattice, but is expensive and has several environmental drawbacks. As most binders include organic compounds, great pains must be taken to process effluent from the production process to ameliorate the negative environmental impact of such compounds. Further, the binder must be cured with an oven, using additional energy and creating additional environmental cleanup costs. While long fibers display fiber to fiber entanglement even without binder, the non-uniformity of the resulting wool packs has long made them commercially undesirable.\nFinally, in addition to the properties of uniformity and integrity, it is desirable for wool packs to exhibit recovery from compression. In the shipping and packaging of insulation products, high compressibility is preferred. It is desirable to compress the wool for shipping and then have it recover rapidly and reliably to the desired size. When the product is compressed, the binder holds firm at fiber to fiber intersections while the glass fibers themselves flex. If the stress upon the fiber increases due to excessive compression, the fiber breaks. Thus, current insulation products are limited in the amount of compression possible while still attaining adequate recovery.\nNonetheless, because long fibers are problematic in nearly all respects, commercial wool insulation products of glass fibers have long used only short straight fibers, despite the various drawbacks of short fibers in lattice non-uniformity, need for binder and related environmental concerns, and limited compressibility. Accordingly, the need remains for further improvements in wool insulation products to improve wool pack properties, reduce cost and eliminate environmental concerns."}
{"text": "For a long time roofing membranes were made from organic felts saturated with bituminous materials such as asphalt or modified bitumen. For several years, those membranes have been improved by adding polyester or fiberglass open grids for strength and fiberglass mats for dimensional stability. See, for example, commonly assigned U.S. Pat. Nos. 4,491,617 and 4,539,254.\nIn this art, it is still generally believed necessary, even with the improved high strength grids, to continue to utilize during manufacture and retain in the final product an organic felt or mat of some kind, such as a polyester mat, to provide adhesion of bitumen to the reinforcing elements, to increase stiffness, to increase bulk, and to provide toughness to resist punctures. In addition, during processing in vats of high temperature bituminous material, such mats hold the grids of the high strength reinforcing elements in place and together inside the bituminous material. For example, without a polyester mat, a high-strength polyester scrim used for reinforcing comes apart in the vat of hot bituminous material during manufacture of the roofing membrane. In addition, high strength grids manufactured without a polyester mat to hold them together have a tendency to \"pull-out\" from the fiberglass mat and come to the surface of the roofing membrane. This occurs while the membrane is being coated with bituminous material or during the cooling stages. It is desirable, however, that the polyester grid remain embedded inside the bitumen.\nAs an example, one current reinforcement for roofing membranes comprises a heavy polyester, open-grid scrim (2.67.times.2.67 yarns per inch, 1000 denier, about 2.0 ounces per square yard) laminated to a lightweight polyester mat (0.5 ounce per square yard, about 17 grams per square meter). This polyester grid/mat product is led with a fiberglass mat (3.75 ounces per square yard) through a vat of bituminous material to create a roofing membrane coated on both sides with bituminous material. During processing the polyester grid is sandwiched between the two mats. One purpose of the polyester mat during this processing is to hold the polyester grid yarns against the fiberglass mat. Without such a lightweight polyester mat, the polyester grid, which is held together with a thermoplastic adhesive (without any mechanical fastening at the cross-overs of the grid), would have a tendency to come apart and/or pull out to the surface of the membrane during impregnation with hot bituminous material.\nA problem with this roofing membrane, however, is delamination of the top and bottom bituminous layers in the final product. While the porosity of the lightweight polyester mat is high and saturates rapidly, the polyester mat does not allow the bituminous layers on either side of it to adhere to each other to a great degree. One method used to reduce this problem is to go to the extra step of needle-punching holes in such mats to promote saturation and adhesion, but this added manufacturing step requires additional expense.\nAccordingly, a need exists to overcome the above noted drawbacks and produce better and more economical products."}
{"text": "The increasing use of cellular telephones and other communication devices in recent years has made the enforcement of federal regulations on their airborne use difficult. Unauthorized use of cellular devices on aircraft poses a risk, for example, of electromagnetic interference (EMI) to the communication components of commercial and cargo aircraft. With the widespread use of cell phones and other communication devices on board aircraft, and the difficulty of enforcing regulations restricting such use, it would be beneficial to automate compliance with the applicable regulations. Automating compliance with the applicable regulations would also aid in power conservation of the devices when they are airborne."}
{"text": "Many electronic devices provide display areas, e.g., computer monitor screens, in which to display data. Often, because of the data's display properties, e.g., its font size, the data is too large to display all at once within the display area. For example, when the data represents a figure, the figure's vertical height or the figure's horizontal width may be, respectively, higher or wider than the vertical height or horizontal width of the display area; or when the data represents text, the number of lines of text may exceed a number of lines that fits within the vertical height of the display area, or the length of each text line may exceed a text length that fits within the horizontal width of the display area.\nIn these instances, to navigate throughout displayable data, an interface, e.g., a graphical user interface (GUI), often provides scroll-bars or toggle-buttons. Scroll-bars are each displayed as a single bar with a predefined length that can be shifted within a predefined space. Often, the predefined space spans across the area in which the data is to be displayed. They may be vertically displayed when the data's vertical height exceeds that of the display area's vertical height, and horizontally displayed when the data's horizontal width exceeds that of the display area's horizontal width. Toggle-buttons are each displayed as a set of two buttons, one placed at a first extremity and another placed at second extremity opposite the first extremity. A user may shift a scroll-bar or depress a toggle-button and thereby indicate a direction in which the data should scroll.\nData applications, especially business data applications, often arrange interrelated data in tables. A table is arranged according to, and illustrates, the structural interrelationship of its data. For example, a table may have columns and rows. Each column may, for example, represent a major category of data. Each row may, for example, represent a sub-category of the major categories. The rows and columns may intersect to form multiple cells. Display of data in a particular cell of the table may indicate that the data relates to the major category of the data's cell's column, and to the sub-category of the of the data's cell's row.\nWhen the sum of all columns' widths exceeds the display area's width or when the sum of all the rows' heights exceeds the display area's height, the data applications provide, respectively, a horizontal scroll-bar or toggle-button to horizontally scroll and/or toggle between the first and last columns, and a vertical scroll-bar or toggle button to vertically scroll and/or toggle between the first and the last rows. For example, a scroll-bar 1 is provided in FIG. 1a. In response to a shift of scroll-bar 1, the entire table, i.e. all columns of the table collectively, may shift. For example, when scroll-bar 1 is shifted to the right, a previously non-displayed column to the right of the column labeled “Document Size” may be displayed; and the previously displayed column labeled “Document ID” may be removed from the display area.\nData, e.g., text, within a row of a particular column may be wider than the particular column, although not necessarily wider than the entire display area. In this instance, not all of the row's text can be viewed at once. For example, in FIG. 1a, the text within each of the cells of the column labeled “Time Stamp” is wider than the column, so that the year is only partially displayed. However, current data applications do not provide for scrolling and/or toggling data within a particular column and/or row, via control of a scroll-bar and/or toggle-buttons. For example, even when scroll-bar 1 is shifted to the right, the text within the “Time Stamp” column will not scroll in relation to the column to reveal the entire year of the cell. Currently, to view all of a cell's text, a user can widen the cell's column; select the text and have a portion of the text superimposed over adjacent columns; or select the text and use an arrow key on a keyboard to move a cursor through and thereby scroll through the data. The first alternative wastes much of the display area. When a column is widened, the column occupies more of the display area, thereby allowing a fewer number of columns to be concurrently displayed. Furthermore, a column may contain numerous rows with data that is not wider than the column's width. For these rows, widening the column width, and thereby using more of the display area is wasteful. The second alternative causes rows of adjacent column to overlap. Furthermore, the second and third alternatives are not as convenient and as easy as scrolling with a scroll-bar or toggle-button. These same alternatives apply when the height of data within the columns of a particular row exceeds the height of the data's particular row."}
{"text": "The present invention relates to phase-locked loops (PLLs) and, in particular, to an integrated PLL having a resistor-less loop filter.\nPLLs are used in integrated circuits, such as application specific integrated circuits (ASICs), for clock synchronization, recovery of serial data streams and frequency synthesis. A typical PLL includes a phase/frequency detector, a charge pump, a loop filter, a voltage-controlled oscillator (VCO) and a frequency divider. The VCO generates a clock signal with a phase and frequency that is a function of the voltage applied to the oscillator. The phase/frequency detector detects a phase and/or frequency difference between the VCO output and the input signal. The phase/frequency detector generates a control signal as a function of the phase difference and applies the control signal to the charge pump which increases or decreases the voltage across the loop filter. This voltage is applied to the VCO for controlling the oscillation frequency and phase of the clock signal.\nA typical first-order loop filter includes a resistor and a capacitor which are coupled together in series between two external pins of the integrated circuit. In a PLL which is fabricated on an integrated circuit, an external loop filter is usually used. An external loop filter is prone to noise pick-up and increases system cost. It is therefore desireable to implement the loop filter on the integrated circuit to form a fully integrated PLL. However, good filter resistors are not readily available in standard digital complementary metal-oxide semiconductor (CMOS) processes. As a result, efforts have been made to eliminate or \"replace\" the filter resistor in on-chip loop filters of fully integrated PLL's. For example, Ilya I. Novof et al. Fully Integrated CMOS Phase-Locked Loop with 15 to 20 MHz Locking Range and +/-50 ps Jitter, IEEE Journal of Solid State Circuits, Vol. 30, No. 11 (November 1995) discloses a phase-locked loop which uses an auxiliary charge pump for replacing the loop filter."}
{"text": "Cultivated crops are affected with various fungal, bacterial and viral diseases as well as a number of insect pests. Many cultivation technical, chemical and biological control methods have been developed in order to control these. The purpose of such methods is to prevent the qualitative and quantitative crop losses caused by plant diseases and other pests.\nIn general the term biological control of plant diseases means the control of plant pathogens by another organism, which is then called a biological control organism of plant diseases, and the preparation made thereof a biological control agent or a biopesticide. The mechanisms of the biological control of plant diseases do vary, and the effect is often based on the cooperative action of many different mechanisms. The control effect may be based on inhibitory chemical agents produced by the organism, sometimes the control organism parasitizes the pathogen or competes with it for growth space and/or the nutrients available.\nThe need of discovering new biological control agents has been increased by the fact that many of the traditional chemical pesticides have turned out to be deleterious to the environment and human beings. A disadvantage of the chemicals is also the fact that many pests have become resistant to one or even a number of pesticides. The development of resistance to biopesticides is instead improbable because the effect thereof is based on a number of mechanisms of different types. The chemicals usually affect faster and more effectively than biopesticides. Biopesticides for their part are often longer-acting than chemicals as their effect is based on a viable and reproducible microorganism.\nThe most important group of biopesticides are bacterial products targeted against insects. Bioinsecticides based on the bacterium Bacillus thuringiensis may be the most commonly used. A biofungicide based on the actinomycete Streptomyces being effective against a number of soil-borne and seed-borne fungal diseases of plants is produced in Finland. Fungi of the genus Trichoderma have also fungicidal activity.\nBacteria of the genus Pseudomonas, especially of the species Pseudomonas fluorescens have been studied a lot and nowadays a great amount of P. fluorescens strains are known which have fungicidal activity. See e.g. published patent applications WO 92/18613, FI 92 1722 and WO 90/01327, as well as EP-patent 228 457.\nIt is known that fungi of the genus Gliocladium have fungicidal activity. Gliocladium virens strains, especially the G. virens strain G1-3, have been described in patent literature. A problem in the use of this species has been, however, that nobody has been able to produce a formulation of this fungus, in which the viability of the fungus would keep on a satisfactory level during the storage of the formulation. In U.S. Pat. Nos. 4,668,512, 4,724,147 and 4,818,530 besides other fungal strains the use of Gliocladium virens strains is described in the preparation of biofungicidal formulations. Also in U.S. Pat. Nos. 5,165,928, 5,194,258 and 5,268,173 the use of Gliocladium virens strains as biofungicides has been described."}
{"text": "Basal-like breast tumors are largely overlapping with the triple-negative subtype defined by the absence of Estrogen Receptor and Progesterone Receptor expression, and Her2 overexpression, and tend to be locally invasive and recur at high frequency (Meyers et al 2011, Pazaiti and Fentiman 2011). Results from animal models and the study of human cancers suggest that E-Cadherin functions as a tumor suppressor and an invasion suppressor (Cowin et al 2005). E-Cadherin is often absent or nonfunctional in invasive breast cancers and this is commonly attributed to epigenetic silencing of the E-Cadherin gene, Cdh1 (Jeanes et al 2008). Therefore, efforts have been directed toward reactivating E-Cadherin expression in cancers (Ou et al 2007, Papageorgis et al 2010). However, data also indicate that E-Cadherin is expressed in some invasive breast cancers where it is mislocalized to cytoplasmic vesicles (Corsino et al 2008, Facina et al 2010). Thus, reinstating E-Cadherin expression may not restore its tumor and invasion suppressive actions unless the posttranslational mechanisms that nullify E-Cadherin function in cancers are also blocked.\nE-Cadherin is a homophilic cell-cell adhesion molecule whose extracellular domain binds to the extracellular domain of E-Cadherin on adjacent cells and whose cytoplasmic domain binds a number of proteins including p120ctn, β-Catenin, and γ-Catenin, also referred to as Junctional Plakoglobin. The latter two proteins bind α-Catenin, which is thought to either directly or indirectly couple the E-Cadherin complex with the Actin cytoskeleton. These multiprotein E-Cadherin-containing complexes play a critical physical role in mediating the formation of adherens junctions in epithelial tissues (Baum and Georgiou 2011), and serve to control the activity of several complex components that also function as transcriptional regulators (Daniel 2007, Heuberger and Birchmeier 2010). As such, mechanisms that converge on E-Cadherin provide an opportunity to address unmet needs in treatment and prevention of disease."}
{"text": "This invention relates generally to the field of fluid jetting systems and relates, more particularly, to the means and methods used to maintain a fluid jet device of a fluid jetting system in proper working order.\nAn example of a fluid jet device with which this invention is concerned is a printer head of an ink jet printer. In some instances, such a printer head relies upon capillary action to move a working fluid (e.g. ink) to the printer head and includes means mounted within the head for directing ink through an orifice toward a target substrate. Such ink-directing means can include an actuator, such a piezoelectric device or an electrostatic membrane, for directing ink through an orifice upon appropriate actuation of the actuator or, in the alternative, can include a thermal device wherein heat which is applied to the ink serves as the mechanism for directing ink through an orifice.\nCommonly, a fluid chamber, or ink flow passageway, is provided in the printer head which conducts ink from a source, by way of a conduit connected between the source and the printer head, and past the ink-directing means to the orifice. During normal operation of the printer head, ink must be present in the ink flow passageway so that operation of the ink-directing means effects a drawing of ink into the passageway and a subsequent pushing of ink, under pressure, through the orifice and toward a target surface. If, however, air enters the ink flow passageway through the orifice (as could be the case if the printing head were accidentally struck or jostled) or if the orifice becomes blocked, for example, by debris or dirt which may become lodged within the orifice, operation of the ink-directing means neither draws additional ink into the passageway nor does it effectively push ink through the orifice. Consequently, for effective operation of the printer head, the ink flow passageway must be devoid of air and the orifices of the printer head must remain free of blockage.\nAir which is present in an ink flow passageway of a printer head and any blockage (surface or internal blockage) of the orifices of a printer head is commonly removed by a purging or head-cleaning operation which requires that additional ink be forced through the conduit and ink flow passageways by way of a purge bulb, pump or other means for forcing ink through the conduit and toward the printer head orifices. Such a purging or head-cleaning process, however, normally pushes ink, as well as air or blockage matter (e.g. debris), through the orifices so that ink which is pushed from the orifices flows downwardly along the front (i.e. the face plate) of the printer head. To prevent the ink which flows downwardly along the front of the printer head from touching or being smeared upon surfaces desired to remain free of ink, the ink is manually wiped from the front of the printing head with an absorbent sheet of material. However, such a purging and subsequent cleaning procedure requires manual intervention in and disruption of the printing operation and is usually a messy, undesirable job. Furthermore, if such a process is required to be performed on a printer head stationed along an assembly line, assembly line production may have to halted in order to satisfactorily service the printer head, thereby causing the loss of production time. It would therefore be desirable to provide a new and improved system and method for maintaining the front of a printer head in a relatively clean condition, even when ink is pushed through the orifices of the printer head during a purging or head-cleaning operation wherein air or blockage material is purged from the ink flow passageways.\nAccordingly, it is an object of the present invention to provide a new and improved system and method for use when purging or cleaning a fluid chamber of a fluid jet apparatus, such as the printer head of an ink jet apparatus, for maintaining the front, or exterior surface, of the fluid jet apparatus in a relatively clean condition.\nAnother object of the present invention is to provide such a system and method which facilitates a purging and cleaning operation in that such operations can be performed routinely upon the fluid jet apparatus without the messiness associated with purging and cleaning operations of the prior art.\nStill another object of the present invention is to provide such a system and method which circumvents the need for wiping the front, or exterior surface, of a fluid jet apparatus during a purging or cleaning operation performed upon the fluid jet apparatus.\nYet another object of the present invention is to provide such a system and method which is well-suited for automatic operation, thereby requiring no manual intervention, and can be performed without disruption of a fluid jetting operation or any assembly line operation with which the system and method are used.\nA further object of the present invention is to provide a system and method which can be used to maintain an orifice of a fluid jet apparatus relatively free of blockage matter, such as debris or dirt.\nA still further object of the present invention is to provide such a system which is uncomplicated in construction, yet effective in operation.\nThis invention resides in an apparatus and method for use when purging air or blockage from a fluid chamber of a fluid jet apparatus during a purging or cleaning operation, wherein the fluid chamber has a chamber wall with an exterior surface, an interior surface adjacent the fluid chamber, and at least one orifice through which fluid is ejected by the fluid jet apparatus during a jetting operation.\nThe apparatus includes a cavity adjacent the chamber wall, wherein the fluid which flows through the at least one orifice during a purging or cleaning operation flows from the at least one orifice into the cavity.\nThe method of the invention includes the steps of providing a cavity adjacent the exterior surface of the chamber wall into which fluid which flows from the at least one orifice during a purging or cleaning operation is permitted to flow from the at least one orifice, and withdrawing fluid which is contained within the cavity.\nIn another aspect of the apparatus and method, the cavity can be used to create a zone of above-atmospheric pressure in a region adjacent the exterior surface of the chamber wall so that leakage of air from said region helps to maintain the at least one orifice in a relatively clean condition. To this end, air is conducted, under pressure, to the interior of the cavity so that any leakage of air from the cavity into the atmosphere effects a flow of air out of the cavity."}
{"text": "1. Technical Field\nThe present disclosure relates to an apparatus for detecting islanding when power supply from a grid power source (commercial power source) is cut off.\n2. Related Art\nRecently, the development of a distributed generation system using solar cells or fuel cells which have a small influence on the environment has been actively conducted in terms of global environment protection. Such a distributed generation system converts DC power generated by the solar cells or the like into AC power with a commercial frequency through an inverter, and supplies the AC power to a system load in connection with a grid power source, while transmitting surplus power to the grid power source.\nThe distributed generation system must prevent reverse charge from the distributed generations to the grid power source when the power supply of the grid power source is stopped. Thus, when the power supply of the grid power source is stopped, the distributed generation system detects islanding of the distributed generations, isolates the distributed generations from the grid power source, and stops operation of the inverter.\nThe method for detecting islanding may include a passive method which detects a rapid change in a voltage waveform or phase when islanding is performed. However, when the balance between output power of the inverter and power consumption of the load is maintained, the passive method cannot detect islanding because no change occurs in the grid when the islanding is performed.\nFurthermore, the method for detecting islanding may include an active method which slightly changes an output voltage or frequency of the inverter and detects that the change becomes prominent when islanding is performed.\nIn the active method, however, when a plurality of distributed power systems are connected in parallel to the grid, the change may be offset to degrade the islanding detection function.\nTechnology for solving such a problem has been disclosed in Japanese Patent Laid-open Publication No. 2012-085384.\nThe technology disclosed in the above-described patent relates to a method for detecting islanding of a distributed generation system which changes DC power of distributed generations into AC power through an inverter and supplies the AC power to a load in connection to a commercial power source of a power grid line. The method controls the inverter to overlap an output current of the inverter at a connection point with the commercial power source with a secondary harmonic current synchronized with a fundamental wave of the output current, and detects islanding based on an even harmonic wave of the fundamental wave included in the grid voltage at the connection point.\nAccording to the method, since the output current of the inverter at the connection point, which is controlled in phase with the grid voltage, overlaps the secondary harmonic current synchronized with the fundamental wave of the output current, the harmonic current overlapped with the output current of the inverter in each of the distributed power systems may be synchronized with the phase of the grid voltage, when the distributed power systems are operated in a state where they are connected in parallel to the grid. Then, since the harmonic currents of the respective distributed power systems do not offset each other, the degradation of the islanding detection function may be prevented.\nIn the above-described technology, however, since both the upper and lower half cycles of the secondary harmonic wave are overlapped, the power quality is degraded.\nFurthermore, while one distributed generation is installed and operated, another distributed generation might need to be additionally installed. In this case, as one of islanding detection methods, the injection directions of reactive power in the plurality of distributed generations may be synchronized through communication. However, when the inverters of the respective distributed generations are manufactured by different makers, the inverters are highly likely to have different control circuits or methods. Thus, the installation of communication lines may become difficult or complex."}
{"text": "With a rapid development of the multimedia society, the technique for semiconductor elements and a display apparatus has got a great and rapid progress.\nIn terms of the displays, an Active Matrix Organic Light-Emitting Diode (AMOLED) display meets those characteristic requirements for a display in the multimedia era, because it has advantages such as no limitations on an angle of view, a low manufacture cost, a high response speed (about hundred times and more of that of the liquid crystal), energy saving, self-luminous, being applicable to a direct-current driving for a portable device, a wider operation temperature scope, and a lighter weight while being miniaturized and thinned along with the hardware devices, and the like. Therefore, the Active Matrix Organic Light-Emitting Diode display has a great potential of development and may be desired to become a next generation of the new flat panel display and to replace the Liquid Crystal Display (LCD).\nCurrently, there are two ways of manufacturing the Active Matrix Organic Light-Emitting Diode display panel, one way of which is to utilize a process technique for Thin Film Transistors (TFTs) of the Low Temperature Poly-Silicon (LTPS), and the other way of which is to utilize a process technique for the Thin Film Transistors (TFTs) of the amorphous silicon (α-Si). Herein, the process technique for the Thin Film Transistors of the Low Temperature Poly-Silicon would need more processes for mask manufacturing craft and thus lead to an increasing of the cost. Therefore, the process technique for Thin Film Transistors of the Low Temperature Poly-Silicon is mostly applied to the panel having a small or medium size, while the process technique for the Thin Film Transistors of the amorphous silicon is mostly applied to the panel having a large size.\nGenerally, for an Active Matrix Organic Light-Emitting Diode display panel manufactured by the process technique for Thin Film Transistors of the Low Temperature Poly-Silicon, the thin film transistors in the pixel circuit may be P-type or N-type, but no matter whether the P-type or N-type transistors are selected to implement the organic light-emitting diode pixel circuit, the current flowing through an organic light-emitting diode may not only change with the changes caused by a long time stress in a turn-on voltage (Voled_th) of the organic light-emitting diode, but also vary with a threshold voltage shift (Vth shift) in the thin film transistor for driving the organic light-emitting diode. As a result, the brightness uniformity and the brightness constancy of the organic light-emitting display may be affected correspondingly."}
{"text": "Transponders are inserted into living beings for the purpose of being able to identify them. A transponder is a transmitter-receiver unit with a memory accommodated in a housing. When the antenna belonging to it is radiated from outside, the stored data are sent out by the transmitter unit, so that remote identification of the animals in question can take place. The stored data are unique for each animal through the coding used. Guns provided with a needle-shaped guide are generally used for the insertion of transponders. An opening is made in the skin in one way or another with said needle-shaped guide, following which the guide is inserted further into the body of the animal in question. The transponder is then ejected by the guide, either by means of a bar or by means of fluid pressure.\nOne method of injection discloses in particular an insertion point for living beings, such as pigs, behind the ear. Although this method of insertion was found to be satisfactory for many animals, it was found that it was unsuitable for animals with moving ears, such as cattle. Unlike pigs, cattle move their ears to a considerable extent, and such animals also have the habit of rubbing their heads along bars of feeder units and the like. These movements cause damage to the transponder and/or force it out of the body. There is also the disadvantage of the transponder being relatively easily accessible."}
{"text": "1. Field of the Invention\nThe present invention concerns a method of degreasing and cleaning an object deposited with oils, as well as an apparatus used therefor. The present invention more particularly concerns such a method and apparatus which can be suitably utilized, for example, as a precleaning in a heat treatment step or as an intermediate cleaning after a hardening step.\n2. Description of the Prior Art\nAs examples of degreasing and cleaning methods which have been applied to metal materials deposited with oils, the following methods are known:\n(1) Alkali cleaning: where cleaning is achieved by dipping the material into a warm aqueous solution, an alcohol solution of sodium hydroxide, or other alkali agent; or by spraying such a solution; alkali cleaning is applied to heavy oil contamination.\n(2) Cleaning with a surface active agent: where cleaning is achieved by dipping the material into a warm aqueous solution of a surface active agent, or by spraying such a solution. Surface active agent cleaning is applied to slight deposition contamination or oil membranes.\n(3) Cleaning by dipping or vapor cleaning with chloro-solvent: a chloro-solvent, such as, for example, 1,1,1-tri-chloroethane, trichloroethylene or perchloroethylene.\n(4) Cleaning with a fluoro-solvent: which involves dipping or vapor cleaning, such as, for example, by using from 113.\nHowever, degreasing and cleaning methods (1) and (2) using the alkali agent or the surface active agent are not generally employed since they involve problems. For example, cleaning performance is poor and quite often causes stains on the surface of an object to be cleaned. Likewise, liquid waste treatment is expensive.\nOn the other hand, degreasing and cleaning methods (3) and (4) using chloro-solvents or fluoro-solvents have high cleaning performance. However, the chloro-solvents involve problems since they show strong toxicity and evaporate greatly. Moreover, the chloro-solvents scatter from cleaning apparatus or treated objects owing to their volatility and thereby contaminate underground water as carcinogenic substances. Accordingly, legal regulations for their use have become severer in recent years. Use of the fluoro-solvents has also been extremely restricted since they form ozone layer destructive substances when released to the atmosphere.\n3. Object of the Invention\nIn view of the above, the present invention has been accomplished taking notice of the foregoing problems in the prior art. Therefore, it is an object hereof to provide a metal degreasing and cleaning method using neither alkali agent nor chloro- or fluoro-solvent which would result in public pollution or circumstantial contamination, as well as an apparatus used for the method."}
{"text": "Lipolytic enzymes such as lipases and phospholipases are used, e.g., in detergents and baking.\nThus, lipases have been used for a number of years as detergent enzymes to remove lipid or fatty stains from clothes and other textiles, particularly a lipase derived from Humicola lanuginosa (EP 258 068 and EP 305 216) sold under the trade name LIPOLASE® (product of NovozymesA/S).\nFatty acid-modified lipases and their use in transesterification have been described. M. Murakami et al., JAOCS, 70 (6), 571–574 (1993); K. Green et al., JAOCS, 75 (11), 1519–1526 (1998).\nIt is also known to add lipases and phospholipases to breadmaking dough. WO 94/04035; WO 98/26057."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a self contained lamp control and lamp package, and relates more specifically to a self contained lamp for automatic light level adjustment.\n2. Description of Related Art\nReferring to FIG. 1, a conventional incandescent light bulb 5 includes base 10, and evacuated envelope or diffuser 14. Diffuser 14, which is typically pear-shaped, surrounds a filament (not shown) that is electrically connected to base 10 for electrical connection to a conventional lamp socket. A well-known base is an Edison screw-base, as shown, which is electrically connectable to an Edison screw socket.\nIn recent years, compact fluorescent lamp packages have been introduced into the market. These lamp packages are particularly desirable because of their energy efficiency.\nReferring to FIG. 2A, an example prior art compact fluorescent lamp package (CFLP) 7 is shown consisting of fluorescent lamp 12, lamp holder 13, and auxiliary housing 15 for interfacing between lamp holder 13 and base 10. Typically, housing 15 also houses a complete ballast 17. However, it should be noted that ballast 17 has also been placed into base 10, while either retaining housing 15 or removing this housing. In most CFLPs an Edison screw-base is used as base 10 for a reliable mechanical and electrical connection to a standard Edison socket, although other sockets are known and used.\nThe typical height of a standard incandescent light bulb 5 is about 4.5 in. However, most conventional CFLPs are typically taller than a conventional incandescent light bulb due to auxiliary elements such as auxiliary housing 15. As a result, a conventional CFLP may extend out of a lampshade or fixture adapted for an incandescent light bulb, causing an undesirable appearance, or the CFLP may not fit within an existing fixture. This extended height, together with the unusual appearance of conventional CFLPs, are believed to adversely affect the acceptance of CFLPs in the marketplace\nAnother issue with the conventional CFLP is that when ballast 17 is housed in auxiliary housing 15, the heat generated by ballast 17 cannot escape very efficiently. Referring, for example, to FIG. 2B, in a typical enclosure or downlight fixture 19, the heat generated from the lamp and ballast of a conventional CFLP 7 collects up in the fixture around the CFLP 7 without any path to escape.\nAs a result, there is a higher ambient temperature around CFLP 7 which causes the internal components of the CFLP 7, especially ballast 17, to run at even higher temperatures, e.g., above 150° C. The increase in the internal heat decreases reliability, causes field failures and limits the use of a conventional CFLP to open rather than enclosed fixtures, all of which further adversely affects the desirability of CFLPs. Nonetheless, it should be noted that when ballast 17 has been placed into base 10 as described above, the base has acted as a heat sink for the ballast, helping to dissipate some of this heat.\nA further issue with the conventional CFLP is that because of its size and operational constraints, such as its high ambient temperature and its use of the conventional Edison socket, it is difficult to realize some desirable features, such as, for example, a dimming feature where the light output of the fluorescent lamp is selected based on a control signal or setting. More specifically, conventional incandescent lights typically use a dimmer control mounted on a wall near an on/off switch for adjusting the light level. These types of dimming controls operate on various principles but all in essence change the characteristic of the electrical power supplied to the light fixture to create various power levels supplied to the lamp, thereby affecting the output lighting level. However, these types of dimming controls are unsuitable for CFLPs because the change in power may extinguish the lamp. As a result, in order to incorporate a dimming feature into the CFLP, a dimmable ballast circuit is typically used where the electronic ballast now includes an additional dimming function that controls the ballast to dim the lamp. However, a mechanism is now needed to control this additional dimming function, thereby making it a difficult challenge to provide a dimming feature to a CFLP.\nFor example, one mechanism that has been used to control a dimmable ballast circuit is to dedicate a wire from a wall-mounted dimmer to the fixture. However, if the CFLP is inserted into an Edison socket, this requires additional wiring. Specifically, Edison screw sockets typically only provide switched utility power on two wires, a line input and a neutral wire. Accordingly, a third wire is now needed to control the dimming function. However, providing a third wire or connection to the Edison screw socket is a highly undesirable additional requirement to realize a dimming feature.\nAnother example mechanism that has been used to control a dimmable ballast circuit is to send control signals from a switch to the fixture directly over the power lines, thereby negating the need for an additional wire. However, this approach is complex, expensive, and space consuming given that the fixture must now include an interface to receive and decode the signal. Nonetheless, it should be noted that some systems that have used this mechanism have also incorporated advantageous features, including the use of an external control module that allows for the remote and automatic control of the lights. For example, a control system may be provided that can control a number of different areas of lighting to automatically change lighting settings depending upon various circumstances, such as evening hours, weekend hours, lengthened days and so forth. However, it should also be noted that these types of lighting systems typically have limited local intelligence (e.g., the fixtures lack sensor feedback that permits localized control) and often control blocks of lights together, rather than individual lights.\nAnother mechanism that has been used to control a dimmable ballast circuit is to include a manual control directly on the casing of the electronic ballast of the CFLP so that dimming can be effected. However, CFLPs designed to replace incandescent bulbs often do not permit easy access to the CFLP or electronic ballast by the user. In addition, once the light level has been set, it is often difficult or inconvenient to readjust.\nA further mechanism that prior fluorescent light systems in general have used to control dimming is to incorporate a self-dimming function that automatically dims the light in reaction to external events. Such self-dimming functions have included the integration of a heat sensor or a photosensitive control element, such as a photocell, with the dimming function. For example, these systems have interconnected a photocell with the dimming function to automatically control the turning on/off of the light and to maintain a constant light level in a room, all based on the changing ambient light levels in the room.\nIn the case of the CFLP in particular, the photocell has been mounted directly on the housing 15. This arrangement creates several advantages as compared to the other dimming mechanisms described above. Specifically, this arrangement allows for individual control of each CFLP and because of the automatic feedback orientation, each CFLP is self-dimming thereby making the CFLP easier and more convenient to use because the user is not required to constantly readjust the light level as conditions change. In addition, when the CFLP is inserted into an Edison socket, the socket does not require additional wiring. This mechanism is also less complex and less expensive than prior mechanisms.\nNonetheless, prior CFLPs that have integrated a sensor, like a photocell, with the dimming function of a dimmable ballast circuit continue to have the same issues as described above with reference to the example CFLP 7. Specifically, these prior CFLPs continue to use a housing 15 to house the dimmable ballast circuit, thereby creating the same height and heat issues described above."}
{"text": "1. Field of the Invention\nThis invention relates to a substrate for superconductive film fabrication to be employed on the occasion of manufacturing oxide superconductive wires through the deposition of a superconductive film on a substrate. This invention also relates to superconductive wires having a superconductive film fabricated on the substrate and to manufacturing methods of the substrate and the superconductive wires.\n2. Description of the Related Art\nConventionally, there have been proposed a many number of ideas for manufacturing oxide superconductive wires through the deposition of a superconductive film on a substrate (see for example, Japanese Patents No. 2614948; No. 3251034; No. 3532253; and No. 3771012; and JP-A 2001-110255 and JP-A 11-3620).\nAmong them, an oxide superconductive wire utilizing an oxide superconductor represented by a composition formula of REBa2Cu3O7-d (wherein RE means rare earth elements, also referred to as 123-type or yttrium-type superconductor), the oxide superconductor being deposited on a tape-like metal substrate to make the wire flexible, is known to attain excellent current characteristics. For this reason, this oxide superconductive wire is one of the oxide superconductive wires that are intensively studied and developed at present. Under the circumstances, the studies on this oxide superconductive wire have been advanced to such a stage that a large number of prototypes concerning electric power equipments and the like utilizing this oxide superconductive wire have been manufactured.\nTherefore, in the art of manufacturing oxide superconductive wires, it is now strongly desired to establish a mass production system aiming at industrialization for the production of superconductive wires, so that we are now faced with a situation where the development of highly reliable and stable process for the production of superconductive wires is urgently needed.\nThe fact that a thin-film oxide superconductor is capable of exhibiting excellent current characteristics has been already confirmed through experiments on the oxide superconductor deposited as a thin film on the surface of a monocryatalline substrate. When a monocryatalline substrate is to be used however, it is impossible to freely bend it and to work it to obtain an elongated strip having a length of several hundreds meters, thus making it unrealistic.\nIn any attempt for the practical use or industrialization of the oxide superconductive wires, it is imperative to work the oxide superconductor into a wire-like member having the same degree of flexibility as metals. For this reason, it has been attempted to deposit an oxide superconductor as a thin film on a tape-like metal substrate to thereby enable the resultant substrate to exhibit a function useful as a practical wire. However, it has been found that, in the case of the superconductive thin film which has fundamentally a polycrystalline structure, it is impossible to obtain excellent current characteristics unless a large number of crystal grains constituting the thin film are uniformly arrayed in the same direction.\nHowever, as conducted in the ion beam assist vapor deposition (IBAD) method, inclined substrate deposition (ISD) method, or a method employing an orientated metal, this problem has been already solved by using a substrate where a large number of crystal grains constituting the substrate are orientated in the same direction as a substrate for superconductive film formation.\nFurther, as a result of considering the matching of lattice constant, reactivity of substrate to the superconductive film, etc. in addition to the orientation of crystal, the development of the superconductive wires is now generally directed to the wires wherein at least one oxide layer called intermediate layer other than a superconductive layer is deposited on a metal substrate and then a superconductive layer is deposited on the intermediate layer, thereby obtaining a wire having a multi-layer structure.\nBy means of this multi-layer structure, it is now possible to obtain wires which are capable of exhibiting excellent current characteristics even on a metal substrate, e.g., a current density of more than 106 A/cm2 at liquid nitrogen temperature (77K) or a current density of more than 107 A/cm2 at liquid helium temperature (4.2K), the wires moreover being flexible.\nHowever, when the practical use or industrialization of the oxide superconductive wires is taken into consideration, it is imperative to further improve the manufacturing speed of the wires. Thus, it is strongly desired, at present, to improve the manufacturing speed of the wires in the art of manufacturing oxide superconductive wires."}
{"text": "Various automated systems for forming and cutting surgical suture tips exist or are known in the art. In one such machine, means for simultaneously advancing in parallel at least six separate strands of suture material, and six independent tensioners for maintaining respective parallel portions of each of the six strands at a preset tension are provided. Once a predetermined length of suture material has been advanced by the advancing means, a horizontal heater bar (positioned perpendicular to the six suture strands) is actuated by an electronically controlled solenoid which moves a planar heater bar into contact with one side of the suture strands for a predetermined dwell time. Once the predetermined dwell time has elapsed, the solenoid retracts the heater bar to its original position, and the heat exposed (or heat-stiffened) section of suture material is advanced to a cutting station. At the cutting station, the heat-stiffened section of suture material is cut at its midpoint, thereby producing a suture with two stiffened ends. Other mechanisms for forming and cutting surgical suture tips are shown in U.S. Pat. Nos. 4,832,025, 4,806,737 and 5,226,336 to Coates. The system described in the Coates patents uses convective or non-contact heating to form suture tips.\nKnown systems for forming and cutting surgical suture tips suffer from several drawbacks. First, such systems typically use heat to stiffen the surgical suture tips. Since the unfinished surgical suture material used by such systems is often coated, the heat applied during the tipping process may melt the coating. Once it has melted, the coating from the unfinished surgical suture material often adheres to the tipping machine, thereby compromising the machine's performance. Another drawback of known systems for forming and cutting surgical suture tips is that such systems typically produce a suture tip, which lacks a substantially uniform cross-section.\nIn addition, such systems are undesirable in that they typically cut the suture tip in an imprecise manner, thereby leaving a cut end, which may be irregular or distorted in shape. From a manufacturing standpoint, suture tips having non-uniform cross-sections and/or irregular or distorted cut ends are undesirable because, among other things, such sutures are difficult to insert into needles. Finally, known systems, which use heat to stiffen surgical suture tips are undesirable because such systems cannot be used with sutures formed from silk.\nIt is therefore an object of the present invention to provide a system for forming surgical suture tips, which system does not use heat in forming the suture tips and which may be used to fuse silk.\nIt is a further object of the present invention to provide a surgical suture having a welded core, which facilitates the easy insertion of the suture tip into a needle.\nIt is a further object of the present invention to provide an automated system and method for manufacturing surgical sutures having tips with welded cores.\nIt is a still further object of the present invention to provide an automated system and method for making surgical sutures with tips having precisely cut ends.\nIt is yet another object of the present invention to provide an apparatus for tipping sutures and method for making suture tips.\nThese and other objects and advantages of the invention will become more fully apparent from the description and claims, which follow or may be learned by the practice of the invention."}
{"text": "This invention relates to a centrifugal fan for sucking air from one of both sides defined along an axis of a revolving shaft of an electric motor and discharging it in a direction perpendicular to an axial direction of the revolving shaft, and more particularly a centrifugal fan exhibiting enhanced waterproof performance.\nA prior art fails to develop a centrifugal fan called a sirocco fan which exhibits increased waterproof performance while being simplified in structure. For example, a centrifugal fan disclosed in Japanese Patent Application Laid-Open Publication No. 148484/1999 is so constructed that a stator which is required to exhibit waterproofness structure is covered with a resin mold.\nHowever, techniques of covering the stator with the resin mold require much time and much cost for formation of the resin mold, to thereby fail to reduce manufacturing costs of the centrifugal fan.\nThe present invention has been made in view of the foregoing disadvantage of the prior art.\nAccordingly, it is an object of the present invention to provide a centrifugal fan which is capable of exhibiting enhanced waterproof performance without using any resin mold while being simplified in structure.\nIt is another object of the present invention to provide a centrifugal fan which is capable of being simplified in structure and exhibiting waterproof performance.\nIt is a further object of the present invention to provide a centrifugal fan which is capable of permitting water collected in a casing to be discharged outwardly of the casing to a degree sufficient not to interfere with operation of the centrifugal fan even when accuracy at which the centrifugal fan is assembled is deteriorated.\nIt is still another object of the present invention to provide a centrifugal fan which is capable of exhibiting increased waterproof performance irrespective of a posture in which the centrifugal fan is mounted or operated.\nIn accordance with the present invention, a centrifugal fan is provided. The centrifugal fan includes an electric motor including a stator having a stator-side magnetic pole and a rotor rotated about a revolving shaft outside the stator. The rotor includes a cup-like member including a bottom wall fixed on the revolving shaft and a peripheral wall arranged so as to extend in an axial direction of the revolving shaft from an outer edge of the bottom wall, and a rotor-side magnetic pole fixed inside the peripheral wall of the cup-like member and arranged opposite to the stator-side magnetic pole of the stator. The centrifugal fan further includes an impeller including a fixed section fixed outside the cup-like member of the rotor of the electric motor, a blade support fixed on the fixed section and arranged so as to extend in a radial direction of the revolving shaft, and a plurality of blades fixed on the blade support so as to suck air on one axial side of the revolving shaft and discharge the air in the radial direction of the revolving shaft; and a casing including a first side wall provided with a suction port which is open to the one axial side of the revolving shaft, a second side wall positioned on the other axial side of the revolving shaft with respect to the first side wall and arranged so as to be spaced from the first side wall at a predetermined interval, and a peripheral wall arranged so as to connect the first side wall and second side wall to each other and form an air passage between the first side wall and the second side wall for guiding air discharged from the impeller to a discharge port which is open in a direction tangential to a direction of rotation of the impeller. The casing receives the electric motor and impeller therein. The term xe2x80x9cone axial sidexe2x80x9d of the revolving shaft used herein is intended to mean one of both sides defined along an axis of the revolving shaft of the motor. Also, the term xe2x80x9cthe other axial sidexe2x80x9d of the revolving shaft used herein is intended to mean the other of both sides defined along the axis of the revolving shaft.\nIn the centrifugal fan generally constructed as described above, the electric motor is arranged in the casing so that the stator is fixed on the second side wall and the cup-like member of the rotor permits the bottom wall to face the suction port of the first side wall. The second side wall of the casing and a wall of the blade support of the impeller have opposite surface portions arranged so as to be opposite to each other in the axial direction of the revolving shaft, respectively. The opposite surface portion of the second side wall of the casing and the opposite surface portion of the wall of the blade support of the impeller each are formed into a configuration which permits the opposite surface portions to cooperate with each other to provide a labyrinth structure which prevents water from intruding into a space defined between the cup-like member and the second wall from an outside in the radial direction of the revolving shaft. The blades are mounted on a surface portion of the blade support opposite to the opposite surface portion of the blade support.\nSuch construction of the present invention permits the centrifugal fan to exhibit satisfactory waterproof performance while being simplified in structure without increasing an overall size thereof and using any resin mold. Also, in the present invention, as described above, the labyrinth structure is arranged at a position at which the blades are mounted, to thereby increase a length of the labyrinth structure and permit a space defined inside the labyrinth structure to be effectively utilized for receiving parts thereon.\nThe labyrinth structure may be constructed in any desired manner. For example, it may be constructed in such a manner that the opposite surface portion of the second side wall is formed with a plurality of cylindrical projections in a manner to project toward the one axial side of the revolving shaft and be arranged concentrically with the revolving shaft, wherein the opposite surface portion of the impeller is integrally formed thereon with a plurality of cylindrical projections in a manner to project toward the other axial side of the revolving shaft, be concentric with the revolving shaft and alternate with the cylindrical projections of the second side walls in the radial direction of the revolving shaft. Such construction permits the cylindrical projections to cooperate with each other to define a gap passage therebetween which provides the labyrinth structure. Also, it facilitates an increase in length of the gap passage. Thus, the centrifugal fan of the present invention exhibits enhanced waterproof performance while being simplified in structure.\nIrrespective of arrangement of the labyrinth structure described above, the centrifugal fan has a likelihood that water is gradually collected in the casing and then intrudes from an inlet of the labyrinth structure through the labyrinth structure into the motor. In order to avoid the problem, at least one of the first and second side walls may be formed at a portion thereof positioned outwardly in the radial direction of the revolving shaft as compared with the labyrinth structure with at least one drainage through-hole for outwardly discharging water collected in the casing. The drainage through-hole may be arranged so as to extend through the side wall in the axial direction. Such arrangement of the drainage through-hole, when water collected in the casing reaches a level of the drainage through-hole, permits the water to be outwardly discharged therethrough from the casing irrespective of a posture in the centrifugal fan is operated and/or operated. This prevents water connected in the casing from intruding through the labyrinth structure into the motor. Water at a level below the drainage through-hole is still remain in the casing. However, the water does not adversely affect operation of the centrifugal fan. The drainage through-hole may be arranged at the peripheral wall. This permits water collected in the casing to be naturally discharged therethrough from the casing. However, when the drainage through-hole of the peripheral wall has a large size, air discharged from the impeller is caused to positively leak therethrough, to thereby deteriorate air feed efficiency and generate noise. Thus, a size of the drainage through-hole and the number thereof are preferably limited to the minimum.\nWhen it is possible to previously know a posture in which the centrifugal fan is operated, a position of the drainage through-hole may be determined depending on the posture. However, in the case of a general-purpose centrifugal fan, a posture in which it is mounted or operated is varied depending on a user. Thus, it is preferable that at least one of the first side wall and second side wall is formed with a plurality of drainage through-holes in a manner to be spaced from each other at intervals in a peripheral direction of the revolving shaft. Such arrangement of the through-holes permits water collected in the casing to be outwardly discharged from the drainage through-holes of the side wall without substantially considering a posture of operation of the centrifugal fan, to thereby enhance general-purpose properties of the centrifugal fan. In particular, in this instance, arrangement of one or more such drainage through-holes at both first side wall and second side wall substantially prevents intrusion of water into the electric motor irrespective of a posture of mounting and/or operation of the centrifugal fan. In this instance, the first side wall and second side wall each may be formed with at least one such drainage through-hole. This prevents intrusion of water into the motor irrespective of a posture of mounting and/or operation of the centrifugal fan.\nFor example, when the drainage through-hole is formed at the side wall of the casing, at least one of the first side wall and second side wall is formed with a plurality of drainage through-holes at a portion thereof positioned in an angular range except a predetermined angular range defined about the revolving shaft and including the discharge port when the first side wall is viewed from the one axial side of the revolving shaft. In this instance, the drainage through-holes may be arranged at portions of the first side wall and second side wall positioned in the predetermined angular range including the discharge port. However, when the centrifugal fan is operated in a posture which may possibly cause water to be collected in the casing under such conditions, water collected in the casing to a degree is permitted to be discharged from the discharge port, so that it is not required to form the drainage through-holes in the angular range. Thus, at least one drainage through-hole may be formed at a position which permits water collected in the casing to be discharged through the drainage through-hole when the centrifugal fan is operated or mounted in a posture which keeps water entering the casing through the suction port from being discharged through the discharge port. Thus, at least one of the first side wall, second side wall and peripheral wall may be formed with at least one drainage through-hole for discharging water in the casing when the centrifugal fan is operated or mounted in a posture which keeps water entering the casing through the suction port from being discharged through the discharge port."}
{"text": "(1) Field of the Invention\nThe invention relates to a door hinge for a motor vehicle with                a column console that can be arranged on a column of a door frame,        a door console that can be attached to a motor vehicle door assigned to the door frame, wherein        the door console and the column console are pin-jointed together via a hinge pin and        the door console or column console is formed from two releasably connected partial units, which rest against each other at assigned contact areas.        \n(2) Prior Art\nDoor hinges for motor vehicles for linking doors or covers, flaps, etc. generally have a door console, which is firmly arranged on a component that is swivel-mounted on a motor vehicle body, and a column console, which is arranged on the vehicle body. The rotating movement of the door and column console relative to each other is normally achieved through a hinge pin, by means of which the consoles are pin-jointed together.\nWithin the optimization of vehicle production and minimization of production costs, the motor vehicle door is meanwhile being removed again after painting after installation on a vehicle body. This takes places in order to facilitate the installation of interior motor vehicle components and the equipping of the motor vehicle door with further hardware components.\nIt is of considerable importance that, after this assembly work is complete, the motor vehicle door is also immediately returned to the original assembly position after painting.\nA motor vehicle door hinge of the initially named type is already known from DE 602 03 143 T2, in which either the door console or the column console has a two-piece structure. These types of motor vehicle door hinges enable the removal of a door or suchlike from the motor vehicle, without separating from each other the components connected together via the hinge pin. After painting, the motor vehicle door hinge is not separated at the hinge pin, but rather via the two-piece door console or column console, wherein the provided fastening means are released. However, known constructions with at least three pieces have the disadvantage that they do not ensure to the required degree the dimensionally accurate restoration of the original installation position."}
{"text": "The present invention relates generally to automated microscopy, and more specifically to improvements in the automatic detection of cell colonies' location on a glass sample slide. Additionally, the glass sample slide may be covered with a coverslip that protects cell colonies from contamination and damage. The coverslip area is that area of a specimen slide where most of focus mapping, image capture and image analysis needs to take place, because cell colonies reside underneath the coverslip. Thus the edges of the coverslip, which denote the area of interest for automated microscopy, need to be reliably detected.\nAt the present, operators typically scan and analyze the entire slide even though colonies of interest may reside only on several isolated spots within the slide. Operators are required to manually identify colonies (by drawing around them) resulting in slow system throughput. Furthermore, focus mapping can be slow or inaccurate due to sparse cell populations on colony slides.\nSome existing methods attempt automated cell analysis of biological specimens by detecting candidate objects. Each slide is first scanned at a low microscope magnification. Candidate objects are identified based on their color, size, and shape; and their location is recorded. The candidate objects are then scanned with higher magnification lens. Thresholding and focusing steps are performed, followed by the morphological processing to identify candidate objects of interest by comparing optical features of the candidate object of interest to a target blob. However, those methods do not use morphological methods that enhance the image of the colonies of interest, neither do they associate the metaphases with the colonies. They also do not disclose a coverslip detection.\nSome other existing methods create a composite image from smaller images. Subsequent image analysis is performed only over the areas of interest within the composite image. Those methods also eliminate the edges that were created by the overlaps or similar imperfections between the subimages caused by mechanical alignment errors. Substantially, those methods could be viewed as bandwidth saving methods. They do not disclose background subtraction, morphological methods for colony detection, thresholding, association of metaphases with the colonies, or the coverslip detection.\nAn accurate identification of the edges of a coverslip on a sample slide continues to be a challenge. Presently, detection methods typically scan and analyze the entire slide, i.e. the areas under and outside of the coverslip, which can be inefficient and time-consuming. Or to reduce scan and analysis time the operators need to accurately place the coverslip in the same position on each slide so that a fixed scan area is applicable to all slides.\nSome methods for detecting a microscope slide coverslip are known. For example, these methods can detect the coverslip by locating all four coverslip edges when those edges satisfy a set of predetermined criteria. However, those methods are rule-based and time consuming, and are not applicable to detecting a coverslip of unknown size and location.\nYet some other methods use non-linear Hough transforms to detect some features of the cell or objects within the cell (e.g., detecting nucleus centre, plasma membrane, etc.). Those methods also use an adjustment of the pixel intensity level to improve feature accuracy, presumably on the suspect edges of the objects of interest. However, those methods detect a presence of the objects within the cell, but not their precise outline, nor do they detect the edges of the coverslip.\nSome other methods detect objects that are similarly shaped using a pre-existing shape library or they detect a grid-like arranged specimens on a slide using Hough transformation. The centroids of the specimen are detected using 2D peak detection algorithms. Column and row orientations are detected followed by the calculation of the overall grid intersection locations. The method can identify the specimens by finding their expected location in the 2D grid. However, those methods do not detect edges of the object (i.e. coverslip edges), neither do they perform any image enhancements, such as, for example, dark field subtraction.\nThere is therefore a need for systems and methods that accurately and automatically detect the location of the coverslip on a microscope slide as well as the location of cell colonies of interest underneath the coverslip."}
{"text": "Advanced forms of information technology (IT) infrastructure, including but not limited to cloud computing platforms, converged infrastructure (CI) platforms, software defined data centers (SDDCs), distributed infrastructure, and other types of processing platforms, are increasing in availability and usage. Management of such IT infrastructure by users, such as customers, companies, businesses, organizations or other enterprises, is a complex task. As IT infrastructure expands and new platforms, software, and other IT resources are released and become available or incorporated into such IT infrastructure, the complex task of IT management is increasingly difficult and costly."}
{"text": "Case hardening is a heat treating practice by which a thin, hardened \"skin\" is formed on a steel article. Case hardening generally designates any of three types of hardening processes: carburizing wherein carbon is added to low carbon steel, nitriding wherein nitrogen is added, and nitrocarburizing wherein both carbon and nitrogen are added to the steel. In carburizing, the article is exposed to a gaseous carbon containing atmosphere (gas carburization) or packed in a carbonaceous powder (pack carburization) or immersed in a cyanide-salt bath (liquid carburizing) and heated so that the carbon diffuses into the surface of the steel to form the skin. In nitriding, the steel article is placed in a heated nitrogen containing atmosphere (gas nitriding), or the article is placed in a heated cyanide bath (liquid nitriding) so that the nitrogen from the gas or bath diffuses into the surface to form the skin. (In both liquid carburizing and liquid nitriding, the medium is a molten cyanide salt bath; the temperature at which the case hardening is carried out determines the result). Ferritic nitro-carburizing (carbonitriding) is a combination of gas carburization and nitriding wherein a nitrogen gas, usually ammonia, and an endothermic gas are added to the carburizing gas.\nThese processes are carried out at temperatures sufficiently high so that nitrogen from the disassociated ammonia or carbon from the carbonaceous powder can diffuse into the steel to form a hardened skin by making iron nitrides and alloys of nitrides and/or carbides through reactions with various trace elements at and just below the surface. Typical gas nitriding temperatures are 925.degree. to 1050.degree. F. whereas typical gas carburization temperatures are above 1600.degree. F. and most often at 1700.degree. F. to 1900.degree. F. Carbonitriding is carried out at temperatures below those required for carburizing, that is at 1400.degree. F. to 1650.degree. F.\nDuring case hardening, anti-nitriding and anti-carburizing compositions are sometimes employed to shield selected metal surfaces from the nitriding or the carburizing medium, whether gas, liquid or powder. These compositions are called \"stop-offs\" or \"resists.\"\nStop-offs are typically employed to prevent a predetermined portion of the steel from forming a hardened skin so as to facilitate later machining operations on the unhardened areas.\nPrevention of case hardening is a matter of preventing nitrogen or carbon from reaching the surface of the steel by providing a physical shield between the steel and the case-hardening atmosphere or environment. The high temperatures at which case hardening processes are carried out present special problems in formulating stop-offs. To be commercially acceptable, a stop-off should exhibit certain desirable performance characteristics: the stop-off must be capable of effective application to the articles at room temperatures (preferably by an easily practiced technique, such as by \"painting\" the article with the stop-off); the stop-off must provide a gas impervious shield which not only resists gaseous penetration at high temperatures but also withstands the stresses and strains caused by thermal expansion and contraction of the underlying article; and the stop-off must be easily removable from the article after the case hardening process is complete.\nCarburizing takes place at substantially higher temperatures than nitriding or ferritic nitrocarburizing so that carburizing stop-off formulations have been different from those used in nitrogen processes.\nOne early stop-off technique was to coat metal with pulverized boric acid. The boric acid melted at the elevated operating temperatures to provide a fluent, adherent glaze-like shield. Such a technique is disclosed in U.S. Pat. No. 1,190,937, which teaches a boric acid coating to prevent the decarbonization of steel during heat tempering. Other carburizing stop-off materials have been proposed which employ various alternative formulations, such as the refractory clay and borax or boric acid formulations in U.S. Pat. Nos. 1,567,632, and 3,151,002. The later patent teaches the use of a synthetic resin lacquer as a binder to hold the boric acid crystals in place until an operating temperature sufficiently elevated to melt the boric acid and to coke the resin is reached.\nNitriding and ferritic nitro-carburizing (nitrogen processes) are carried out at temperatures lower than those for carburizing. Nitrogen process frequently employ a layer of tin on the surfaces of the article to be shielded. Electroplating was one technique of providing a tin layer, and where electroplating was not feasible, various tin-bearing coatings have been employed to shield parts of the article. When the article was brought to an elevated operating temperature, the tin in the coating melted to form a coherent, fluent shield. A tin technique, albeit for prevention of carburization in a cyanide bath, is disclosed in U.S. Pat. No. 2,485,176, which teaches a copper/tin (bronze) composition in an organic, vehicle.\nIn order to be easily applied by spraying or painting, some prior art stop-offs were formulated with combustible or flammable resins, as taught in U.S. Pat. No. 2,485,176, wherein the finely powdered metal particles are dispersed in an organic vehicle so as to be sprayable or paintable on surfaces not required to be case-hardened. If further thinning or dilution of these stop-offs was needed, for example after a period of storage or standing open so as to cause the resin to evaporate, additional flammable resin or solvent therefor had to be handled. Thus, use of many prior stop-offs involved the risk of fire or explosion. It can be appreciated that such stop-offs required special ventilation of the work place as well as special handling and storage facilities. Lastly, flammable stop-offs present transportation problems in that they require special packaging, labeling and handling.\nIn order to ameliorate the dangers and difficulties attendant the use of flammable stop-offs, various proposals have been made for providing non-flammable formulations. U.S. Pat. No. 3,178,321 teaches a wide temperature range (i.e., from about 800.degree. F. to about 1900.degree. F.) protective coating comprising refractory clays in a water-dispersible resin. U.S. Pat. No. 3,454,433 discloses low temperature (i.e., below about 1600.degree. F.) protective ceramic coatings comprising frits and refractory materials in both organic solvent type binder carrier and water-base binder carrier. U.S. Pat. No. 3,661,820 teaches an anti-carburization compound containing boric acid or borax in \"water reducible\" resins. Although the greater water content of these stop-offs may have improved safety, there was a trade-off in terms of drying techniques and times. The '321 patent reports oven drying at 180.degree. F. In the '433 patent, air drying is followed by baking the water-soluble resinous varnish-containing formulation at about 180.degree. C. to thermoset the binder; drying times are reported in the '820 patent to be around eight hours at room or factory temperatures. Satisfactory use of such stop-offs required either the provision of special elevated temperature drying rooms to hasten drying or a great deal of advance planning coupled with storage space so that the coated articles awaiting treatment could be stored away during drying so as to avoid damage to the wet stop-off surfaces.\nAnother disadvantage encountered with prior formulations for water-based stop-offs included short-shelf life in contrast to organic solvent based stop-offs; problems with prior formulations included such settling during storage that the particulates could not be stirred up again, as well as rapid hardening after opening the container.\nIt would be advantageous, therefore, to have a truly water-based stop-off useful in nitrogen processes and characterized by ease of application and quickness of drying before heat treatment, ease of removal afterward, and enhanced safety in the work place."}
{"text": "The invention relates to a stand for vertically receiving rod-shaped materials, comprising a base plate that has at least two diametrically opposed clamping jaws which are disposed on a substantially horizontal plane and are used for receiving a rod. At least one of said jaws can be displaced radially in an outward direction counter to the force of a spring.\nA device of this generic concept is known from DE-A-36 23 693, which discloses at least two diametrically opposed clamping jaws that grasp the vertical rod, for example, a Christmas tree or a flagpole, whereby at least one of the clamping jaws is spring-loaded and can be displaced radially in an outward direction. As far as practical handling is concerned, it proves difficult to open the displaceable clamping jaw counter to the force of the spring by hand and, at the same time, using the other hand, to introduce the rod between the opened clamping jaws, especially if the object is a flagpole, a Christmas tree or a larger mast. The help of a second person is usually required, whoxe2x80x94after the first person has opened the clamping jawsxe2x80x94positions the rod or the Christmas tree between the clamping jaws.\nOn this basis, the invention has the object of providing a stand which can be used effortlessly by a single person and which brings a rod into a secure vertical position and fixes it there.\nIn accordance with the invention, this task is solved therein that the displaceable clamping jaw is pivotally connected to one end of a lever, whereby the opposite end thereof is guided and displaced in an abutment in a substantially vertical manner and is located above the horizontal plane defined by the direction of movement of the clamping jaw.\nThe invention\"\"s central idea consists therein to displace the clamping jaw to be actuated for opening via a lever, which is disposed in a special way and manner and is guided displaceably in an abutment. Here, the lever, most simply conceivable as a straight bar, is pivotally connected directly or indirectly to the displaceable clamping jaw, whereby the opposite end of the lever is guided in an abutment that facilitates substantially vertical displacement. The arrangement of the different possible positions of the end of the lever in the abutment is to be effected so that said end is located above the horizontal plane defined by the direction of movement of the clamping jaw during every phase of motion. Consequently, a downward movement of the end guided in the abutment presses the clamping jaw away from the abutment, thereby opening the stand or releasing the rod. Within the meaning of the invention, the term xe2x80x98vertically displaceablexe2x80x99 denotes that the abutment may comprise a strictly vertical guidance, but also a guidance with a specific inclination in relation to the perpendicular. The displacement of the end of the lever within the abutment can also be achieved by a downward movement of the foot, that is, the operator steps on the end of the lever, whereby both hands are free to insert or remove the rod from the clamping jaws. To realise the inventive idea, it is sufficient to employ two clamping jaws, one of which is displaceable. Moreover, a larger number of clamping jaws can be provided, whereby one or also several are disposed displaceably and displaced synchronously.\nAs regards the use of the application force produced by the springs, it is optimum if rigid and displaceable clamping jaws are disposed diametrically opposite each other. The application force is fully compensated by the rigid, opposing clamping jaw in the sense of a bearing. One then obtains a optimum relationship of forces for configuring the clamping jaws.\nIn an advantageous embodiment, the lever is configured like a frame, whereby one lateral length corresponds to the end of the lever connected pivotally to the displaceable clamping jaw and the opposing lateral length produces the connection to the abutment. The parts of the frame located in between, which extend essentially perpendicularly thereto, give a closed form that circumscribes the reception of the rod.\nTo improve the guidance, and for lateral support, in a further embodiment, the lever (or also the frame) is guided on both sides in the region of the swivelling axis in a guiding groove in the base plate by a cross brace, which is aligned in the direction of movement of the clamping jaw. Although further possibilities for guiding the displaceable clamping jaw are conceivable, this comprises a preferred option because it provides an optimum manner of support.\nThe constructional, concrete form of the abutment is configured as an essentially vertically aligned inclination, along which the end of the lever slides. By means of the force of the spring upon the displaceable clamping jaw, and the opposite end of the lever connected thereto, the end of the lever on the abutment side is pressed against the inclination. When a downward movement occurs, the opposite clamping jaw is pressed away from the abutment, whereby the downward movement can be achieved by stepping thereupon. To achieve the best transmission of forces, with as little friction as possible, the end of the lever rests against the inclination via a roller so that a loss of force is minimised by means of rolling friction.\nTo facilitate the downward movement of the lever in the region of the abutment, a tread surface is disposed there, upon which one can step and thus the downward movement of the lever can be effected comfortably by means of application of force, that is, by moving one\"\"s foot downward or through the operator\"\"s weight.\nTo make it possible to place both feet thereupon, and to enable the application of one\"\"s total weight while avoiding overturning moments, it is proposed to configure two tread surfaces symmetrically in relation to the abutment.\nFor utilising the inventive stand, the shape of the rod cross section is basically optional. To produce optimum, that is, large-surface contact between the rod and the clamping jaws, it is advantageous to select the interior shape of the clamping jaws so that it is commensurate with the shape of the cross section of the rod to be inserted. In the case of a round rod, clamping jaws that have a circular form with the corresponding radius are especially used. This embodiment is also recommended if the material used for the rod is sensitive and should be largely free of pressure marks.\nWith other materials, the option is given to select a rhombus shape for the closed clamping jaws. Then, the contact between the clamping jaw and the rod does not have to be produced across the whole circumference so that this shape proves advisable especially if the shape of the rod cross section is irregular.\nIf soft materials are used, support and force transmission can be improved by integrating a spike inside the clamping jaws, that is, on the side facing the rod. When the clamping jaws are closed, the spike digs into soft materials, for example, wood, and in this way establishes intimate contact."}
{"text": "The alkaloids obtainable from Vinca rosea represent one of the most productive areas of chemistry for drugs which adversely affect the growth of experimental malignancies in mammals. Initially, only some of the alkaloids, obtainable from the leaves of the plant by extraction and purifiable by chromatography, were found to be active. These active antineoplastic alkaloids obtained directly from the leaves of the vinca plant incude VLB (Vinblastine, vincaleucoblastine), vincristine (leurocristine), leurosine (vinleurosine), leurosidine (vinrosidine), leuroformine (formylleurosine) and deoxy VLB \"A\" and \"B\" (4'-deoxy VLB and 4'-deoxyleurosidine).\nChemical modification of the Vinca alkaloids started slowly for several reasons. In the first place, the molecular structures involved are extremely complex, and chemists were slow to find reactions which modified one specific functional group of the molecule without affecting other groups. Secondly, dimeric alkaloids lacking desirable chemotherapeutic properties had been recovered or produced from Vinca rosea extracts, and a determination of their structures had led to the conclusion that these inactive compounds were closely related structurally to, and even isomeric with, one or more of the active alkaloids. Thus, it appeared that small chemical changes in the known anticancer alkaloids could have a profound effect on antineoplastic activity.\nBecause of these restrictions, modification of the indole-dihydroindole alkaloids obtained from Vinca rosea has centered around only three areas of the molecule: C-3, C-4' and C-4. Considering C-3 modification first, one of the more recent, and more successful, modifications of the basic indole-dihydroindole structure has been the preparation of C-3 carboxamide and carboxhydrazide derivatives, most of which turned out to be active anti-tumor agents. [See U.S. Pat. No. 4,166,810, and Conrad et al. J. Med. Chem., 22, 391 (1979)]. 4-Desacetyl VLB 3-carboxamide (vindesine) is currently being marketed in several European countries as an oncolytic agent. It is said to be effective in treating some vincristine-resistant leukemias in addition to many common neoplasms including germ-cell tumors. Reaction of the 3-hydroxy and 3-ester functions with an isocyanate has produced the corresponding oxazolidinedione derivatives, one of which, the N-chloroethyl derivative--vinzolidine--is currently undergoing a clinical trial. These oxazolidinedione derivatives are disclosed in Miller and Gutowski, RE 30,560, reissued Mar. 31, 1981.\nA second area of the molecule which has been modified is the C-4' functionality. A majority of these modifications have been based on the 3',4'-anhydro derivative, makeable both by coupling vindoline and catharanthine via a modified Polonovski reaction--Potier et al. J.C.S. Chem. Comm., 670, (1975)--and by dehydrating VLB or leurosidine--Gutowski and Miller, U.S. Pat. No. 4,029,663. The dehydration reaction produces two exo-double bond isomers in addition to the .DELTA..sup.3',4' -anhydro derivative. Functionalization of any one of these double bonds to form epoxides, diols, etc. has constituted the chief chemical modifications undertaken at C-4'.\nThe third region of the indole-dihydroindole which has been modified successfully is C-4. In the first place, hydrolysis of the acetoxy group, present in all the above vinca alkaloids, yields active antineoplastic 4-desacetyl derivatives. (Vindesine, a C-3 carboxamide, is a 4-desacetyl derivative.) Secondly, Hargrove, U.S. Pat. Nos. 3,387,001 and 3,392,173 prepared novel 4-acyl derivatives of 4-desacetyl VLB, 4-desacetyl vincristine, etc. Among these new derivatives was 4-chloroacetyl VLB, which compound could be reacted with amines, for example, dimethylamine, to yield a potent anticancer drug, vinglycinate, N,N-di-methyl 4-glycinyl VLB. In a different modification, Wright and Neuss, U.S. Pat. No. 4,122,082, oxidized the 4-hydroxyl of 4-desacetyl VLB to a 4-keto compound, and Thompson, U.S. Pat. No. 4,195,022, reduced this ketone to the 4-epihydroxy (4.alpha.-hydroxy) derivative, also a compound with anticancer activity.\nIndole-dihydroindole bridged dimers; i.e., the same or different alkaloid moieties bridged thru the 3-carboxyl via a bis-amide are described in Conrad and Gerzon, U.S. Pat. No. 4,199,504. Otherwise, indole-dihydroindole vinca alkaloid dimers have not been bridged through other positions in the molecule to form vinca tetramers.\nVLB and vincristine have been conjugated with proteins to form antigens useful in radioimmune assays. 4-Desacetyl VLB 3-carboxazide (desacetyl vinblastinoic azide) and the corresponding vincristine compound have been the derivatives employed; see Conrad et al., J. Med. Chem., 22, 391 (1979), European Pat. No. 41,935, Abstract 182, FACSS, Oct. 6, 1975, and U.S. Pat. No. 4,203,898 for illustrations of this reaction."}
{"text": "Minimally invasive surgery (MIS) is the performance of surgery through incisions that are considerably smaller than incisions used in traditional surgical approaches. For example, in an orthopedic application such as total knee replacement surgery, an MIS incision length may be in a range of about 4 to 6 inches, whereas an incision length in traditional total knee surgery is typically in a range of about 6 to 12 inches. As a result of the smaller incision length, MIS procedures are generally less invasive than traditional surgical approaches, which minimizes trauma to soft tissue, reduces post-operative pain, promotes earlier mobilization, shortens hospital stays, and speeds rehabilitation.\nMIS presents several challenges for a surgeon. For example, in minimally invasive orthopedic joint replacement, the small incision size may reduce the surgeon's ability to view and access the anatomy, which may increase the complexity of sculpting bone and assessing proper implant position. As a result, accurate placement of implants may be difficult. Conventional techniques for counteracting these problems include, for example, surgical navigation, positioning the subject patient limb for optimal joint exposure, and employing specially designed, downsized instrumentation and complex surgical techniques. Such techniques, however, typically require a large amount of specialized instrumentation, a lengthy training process, and a high degree of skill. Moreover, operative results for a single surgeon and among various surgeons are not sufficiently predictable, repeatable, and/or accurate. As a result, implant performance and longevity varies among patients.\nTo assist with MIS and conventional surgical techniques, advancements have been made in surgical instrumentation, and in technologies for understanding the spatial and rotational relationships between surgical instruments and tissue structures with which they are intervening during surgery. For example, instruments for calcified tissue intervention that are smaller, lighter, and more maneuverable than conventional instruments have become available, such as handheld instruments configured to be substantially or wholly supported manually as an operator creates one or more holes, contours, etc. in a subject bony tissue structure. In certain surgical scenarios, it is desirable to be able to use such handheld type instrumentation while understanding where the working end of the pertinent tools are relative to the anatomy. In particular, it is desirable to be able to control the intervention such that there are no aberrant aspects, wherein bone or other tissue is removed outside of the surgical plan, as in a situation wherein a surgical operator has a hand tremor that mistakenly takes the cutting instrument off path, or wherein a patient moves unexpectedly, thereby taking the instrument off path relative to subject tissue structure. There is a need for handheld systems that are capable of assisting a surgeon or other operator intraoperatively by compensating for aberrant movements or changes pertinent to the spatial relationship between associated instruments and tissue structures."}
{"text": "The disclosures, including the claims, figures and/or drawings, of each and every patent, patent application, and publication cited herein are hereby incorporated herein by reference in their entireties.\nThe following description includes information that may be useful in understanding the present invention. It is not an admission that any of the information provided herein is prior art or relevant to the presently claimed inventions, or that any publication specifically or implicitly referenced is prior art.\nThere are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The goal is to combine in a single variety or hybrid an improved combination of desirable traits from the parental germplasm. These important traits may include higher yield, resistance to diseases and insects, better stalks and roots, tolerance to drought and heat, reduction of grain moisture at harvest as well as better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination and stand establishment, growth rate, maturity and plant and ear height is important.\nChoice of breeding or selection methods depends on the mode of plant reproduction, the heritability of the trait(s) being improved, and the type of cultivar used commercially (e.g., F1 hybrid cultivar, pureline cultivar, etc.). For highly heritable traits, a choice of superior individual plants evaluated at a single location will be effective, whereas for traits with low heritability, selection should be based on mean values obtained from replicated evaluations of families of related plants. Popular selection methods commonly include pedigree selection, modified pedigree selection, mass selection, recurrent selection, and backcross breeding.\nThe complexity of inheritance influences choice of breeding method. Backcross breeding is used to transfer one or a few favorable genes for a heritable trait into a desirable cultivar. This approach has been used extensively for breeding disease-resistant cultivars; nevertheless, it is also suitable for the adjustment and selection of morphological characters, color characteristics and simply inherited quantitative characters such as earliness, plant height or seed size and shape. Various recurrent selection techniques are used to improve quantitatively inherited traits controlled by numerous genes. The use of recurrent selection in self-pollinating crops depends on the ease of pollination, the frequency of successful hybrids from each pollination, and the number of hybrid offspring from each successful cross.\nEach breeding program should include a periodic, objective evaluation of the efficiency of the breeding procedure. Evaluation criteria vary depending on the goal and objectives, but should include gain from selection per year based on comparisons to an appropriate standard, overall value of the advanced breeding lines, and number of successful cultivars produced per unit of input (e.g., per year, per dollar expended, etc.).\nPromising advanced breeding lines are thoroughly tested per se and in hybrid combination and compared to appropriate standards in environments representative of the commercial target area(s) for three or more years. The best lines are candidates for use as parents in new commercial cultivars; those still deficient in a few traits may be used as parents to produce new populations for further selection.\nThese processes, which lead to the final step of marketing and distribution, usually take from eight to twelve years from the time the first cross is made. Therefore, development of new cultivars is a time-consuming process that requires precise forward planning, efficient use of resources, and a focus on clear objectives.\nA most difficult task is the identification of individuals that are genetically superior, because for most traits the true genotypic value is masked by other confounding plant traits or environmental factors. One method of identifying a superior plant is to observe its performance relative to other experimental plants and to a widely grown standard cultivar. If a single observation is inconclusive, replicated observations provide a better estimate of its genetic worth.\nThe goal of corn breeding is to develop new, unique and superior corn inbred lines and hybrids. The breeder initially selects and crosses two or more parental lines, followed by repeated self pollination or selfing and selection, producing many new genetic combinations. Another method used to develop new, unique and superior corn inbred lines and hybrids occurs when the breeder selects and crosses two or more parental lines, followed by haploid induction and chromosome doubling that results in the development of dihaploid inbred lines. The breeder can theoretically generate billions of different genetic combinations via crossing, selfing and mutations and the same is true for the utilization of the dihaploid breeding method.\nEach year, the plant breeder selects the germplasm to advance to the next generation. This germplasm is grown under unique and different geographical, climatic and soil conditions, and further selections are then made, during and at the end of the growing season. The inbred lines which are developed are unpredictable. This unpredictability is because the breeder's selection occurs in unique environments, with no control at the DNA level (using conventional breeding procedures or dihaploid breeding procedures), and with millions of different possible genetic combinations being generated. A breeder of ordinary skill in the art cannot predict the final resulting lines he develops, except possibly in a very gross and general fashion. This unpredictability results in the expenditure of large research funds to develop a superior new corn inbred line.\nThe development of commercial corn hybrids requires the development of homozygous inbred lines, the crossing of these lines, and the evaluation of the crosses. Pedigree breeding and recurrent selection breeding methods are used to develop inbred lines from breeding populations. Breeding programs combine desirable traits from two or more inbred lines or various broad-based sources into breeding pools from which inbred lines are developed by selfing and selection of desired phenotypes or through the dihaploid breeding method followed by the selection of desired phenotypes. The new inbreds are crossed with other inbred lines and the hybrids from these crosses are evaluated to determine which have commercial potential.\nPedigree breeding is used commonly for the improvement of self-pollinating crops or inbred lines of cross-pollinating crops. Two parents which possess favorable, complementary traits are crossed to produce an F1. An F2 population is produced by selfing one or several F1s or by intercrossing two F1s (sib mating). Selection of the best individuals is usually begun in the F2 population; then, beginning in the F3, the best individuals in the best families are selected. Replicated testing of families, or hybrid combinations involving individuals of these families, often follows in the F4 generation to improve the effectiveness of selection for traits with low heritability. At an advanced stage of inbreeding (i.e., F6 and F7), the best lines or mixtures of phenotypically similar lines are tested for potential release as new cultivars. Similarly, the development of new inbred lines through the dihaploid system requires the selection of the best inbreds followed by four to five years of testing in hybrid combinations in replicated plots.\nMass and recurrent selections can be used to improve populations of either self- or cross-pollinating crops. A genetically variable population of heterozygous individuals is either identified or created by intercrossing several different parents. The best plants are selected based on individual superiority, outstanding progeny, or excellent combining ability. The selected plants are intercrossed to produce a new population in which further cycles of selection are continued.\nBackcross breeding has been used to transfer genes for a simply inherited, highly heritable trait into a desirable cultivar or inbred line which is the recurrent parent. The source of the trait to be transferred is called the donor parent. The resulting plant is expected to have the attributes of the recurrent parent (e.g., cultivar) and the desirable trait transferred from the donor parent. After the initial cross, individuals possessing the phenotype of the donor parent are selected and repeatedly crossed (backcrossed) to the recurrent parent. The resulting plant is expected to have the attributes of the recurrent parent (e.g., cultivar) and the desirable trait transferred from the donor parent.\nDescriptions of other breeding methods that are commonly used for different traits and crops can be found in one of several reference books (e.g., R. W. Allard, 1960, Principles of Plant Breeding, John Wiley and Son; N. W. Simmonds, 1979, Principles of Crop Improvement, Longman Group Limited; W. R. Fehr, 1987, Principles of Crop Development, Macmillan Publishing Co.; N. F. Jensen, 1988, Plant Breeding Methodology, John Wiley & Sons).\nProper testing should detect any major faults and establish the level of superiority or improvement over current cultivars. In addition to showing superior performance, there must be a demand for a new cultivar that is compatible with industry standards or which creates a new market. The introduction of a new cultivar will incur additional costs to the seed producer, the grower, processor and consumer for special advertising and marketing, altered seed and commercial production practices, and new product utilization. The testing preceding release of a new cultivar should take into consideration research and development costs as well as technical superiority of the final cultivar. For seed-propagated cultivars, it must be feasible to produce seed easily and economically.\nOnce the inbreds that give the best hybrid performance have been identified, the hybrid seed can be reproduced indefinitely as long as the homogeneity of the inbred parent is maintained. A single-cross hybrid is produced when two inbred lines are crossed to produce the F1 progeny. A double-cross hybrid is produced from four inbred lines crossed in pairs (A×B and C×D) and then the two F1 hybrids are crossed again (A×B)×(C×D). Much of the hybrid vigor exhibited by F1 hybrids is lost in the next generation (F2). Consequently, seed from hybrid varieties is not used for planting stock.\nHybrid corn seed is typically produced by a male sterility system or by incorporating manual or mechanical detasseling. Alternate strips of two corn inbreds are planted in a field, and the pollen-bearing tassels are removed from one of the inbreds (female). Providing that there is sufficient isolation from sources of foreign corn pollen, the ears of the detasseled inbred will be fertilized only from the other inbred (male), and the resulting seed is therefore hybrid and will form hybrid plants.\nThe laborious, and occasionally unreliable, detasseling process can be avoided by using cytoplasmic male-sterile (CMS) inbreds. Plants of a CMS inbred are male sterile as a result of factors resulting from the cytoplasmic, as opposed to the nuclear, genome. Thus, this characteristic is inherited exclusively through the female parent in corn plants, since only the female provides cytoplasm to the fertilized seed. CMS plants are fertilized with pollen from another inbred that is not male-sterile. Pollen from the second inbred may or may not contribute genes that make the hybrid plants male-fertile. Seed from detasseled fertile corn and CMS produced seed of the same hybrid can be blended to insure that adequate pollen loads are available for fertilization when the hybrid plants are grown.\nThere are several methods of conferring genetic male sterility available, such as multiple mutant genes at separate locations within the genome that confer male sterility, as disclosed in U.S. Pat. Nos. 4,654,465 and 4,727,219 to Brar et al. These and all patents referred to are incorporated by reference. In addition to these methods, Albertsen et al., U.S. Pat. No. 5,432,068, have developed a system of nuclear male sterility which includes: identifying a gene which is critical to male fertility, silencing this native gene which is critical to male fertility; removing the native promoter from the essential male fertility gene and replacing it with an inducible promoter; inserting this genetically engineered gene back into the plant; and thus creating a plant that is male sterile because the inducible promoter is not “on” resulting in the male fertility gene not being transcribed. Fertility is restored by inducing, or turning “on,” the promoter, which in turn allows the gene that confers male fertility to be transcribed.\nThere are many other methods of conferring genetic male sterility in the art, each with its own benefits and drawbacks. These methods use a variety of approaches such as delivering into the plant a gene encoding a cytotoxic substance associated with a male tissue specific promoter or an anti-sense system in which a gene critical to fertility is identified and an antisense to that gene is inserted in the plant (see, Fabinjanski, et al. EPO 89/0301053.8 publication number 329,308 and PCT application PCT/CA90/00037 published as WO 90/08828).\nAnother version useful in controlling male sterility makes use of gametocides. Gametocides are not a genetic system, but rather a topical application of chemicals. These chemicals affect cells that are critical to male fertility. The application of these chemicals affects fertility in the plants only for the growing season in which the gametocide is applied (see Carlson, G. R., U.S. Pat. No. 4,936,904). Application of the gametocide, timing of the application, and genotype often limit the usefulness of the approach.\nCorn is an important and valuable field crop. Thus, a continuing goal of plant breeders is to develop stable, high yielding corn hybrids that are agronomically sound. The reasons for this goal are obviously to maximize the amount of ears or kernels produced on the land used and to supply food for both humans and animals. To accomplish this goal, the corn breeder must select and develop corn plants that have the traits that result in superior parental lines for producing hybrids.\nThe foregoing examples of the related art and limitations related therewith are intended to be illustrative and not exclusive. Other limitations of the related art will become apparent to those of skill in the art upon a reading of the specification."}
{"text": "The present invention relates to video camera units and protective enclosures for video cameras, and in particular, to a video camera unit especially adapted for use in vehicles, for example, school buses. The invention also relates to a power circuit and controller circuit for portable video cameras for use in vehicles.\nVideo camera units have increasingly been used in vehicles, especially school buses, in order to provide surveillance of the students being transported in the school bus. By using a video camera unit in a school bus, unruly students can be supervised, and it has been found that the mere presences of a video camera unit, or even a non-functional empty enclosure or mock-up of such a unit, is effective in maintaining order on school bus vehicles."}
{"text": "This invention relates to a method for controlling wet-process flue gas desulfurization, and more particularly to a method for controlling wet-process flue gas desulfurization suitable for optimum operating control of a whole plant including a combustion unit such as a boiler, etc. and a desulfurization unit."}
{"text": "The present disclosure relates to a semiconductor structure, and particularly to a semiconductor structure including locally thinned semiconductor fins and a method for manufacturing the same.\nFin field effect transistors are widely employed in advanced semiconductor circuits for their superior performance over planar field effect transistors. Fin field effect transistors provide a high on-current per area and full depletion of a channel during operation. Typically, the width of the semiconductor fins employed for the field effect transistors is the minimum lithographically printable dimension, which is referred to as the critical dimension. Lithographic bias and etch bias that are introduced during lithographic patterning and pattern transfer can cause variations in the fin widths between semiconductor fins in a nested environment and semiconductor fins in an isolated environment. Thus, a method of selectively adjusting the widths of a group of semiconductor fins without affecting other semiconductor fins is desired.\nFurther, typical semiconductor processing flows provide semiconductor fins of identical height. This feature tends to quantize the on-current of fin field effect transistors. While the on-current of fin field effect transistors can be increased in multiples of the on-current of a fin field effect transistor employing a single fin, providing a fin field effect transistor having an on-current at a fractional multiple of the on-current of a fin field effect transistor employing a single fin is a challenge for conventional semiconductor processing sequences. Thus, a method is desired for providing field effect transistors having an on-current that is a non-integer multiple of the on-current of a fin field effect transistor employing a single fin."}
{"text": "Semiconductor and other types of electronic devices are often encapsulated wholly or partly in plastic resin to provide environmental protection and facilitate external connection to the devices. For convenience of explanation and not intended to be limiting, the present invention is described for semiconductor devices, but persons of skill in the art will understand that the present invention applies to any type of electronic device that is substantially in chip form. Accordingly, such other types of devices including the non-limiting examples given below, are intended to be included in the terms “device”, “electronic device”, “semiconductor device” and “integrated circuit” whether singular or plural, and the terms “device”, “die” and “chip” are intended to be substantially equivalent. Non-limiting examples of suitable devices are semiconductor integrated circuits, individual semiconductor devices, piezoelectric devices, magnetostrictive devices, solid state filters, magnetic tunneling structures, integrated passive devices such as capacitors, resistors and inductors, and combinations and arrays of any and all of these types of devices and elements. Further, the present invention does not depend upon the types of die or chips being used nor the materials of which they are constructed provided that such materials withstand the encapsulation process.\nIn certain types of electronic device packaging where connections to multiple devices included in the package are made after encapsulation, there is a problem referred to as warping that can occur during encapsulation. Warping is of particular concern in electronic assemblies that are in the form of a comparatively flat or planar panel whose device electrical connections are exposed on a principal surface. It is often desired to form an integrated electronic assembly by interconnecting the various devices in the panel using planar processing technology. If the panel has warped during encapsulation, the process of adding the interconnects becomes difficult, which can affect overall yield and cost. Thus, control or elimination of warping is important to achieving high manufacturing yields and low manufacturing costs in such encapsulated planar assemblies.\nAccordingly, it is desirable to provide packaging for electronic devices that avoids or mitigates the adverse effects of warping during encapsulation. It is further desirable that the packaging is suitable for use with arrays containing multiple devices and/or multiple types of devices and especially device arrays where it is desired that the primary faces of the devices are available for electrical connections thereto by planar processing or the like after the devices are fixed in the encapsulation. In addition, it is desirable that the methods, materials and structures employed be compatible with available manufacturing capabilities and materials and not require substantial modifications of manufacturing procedures or substantially increase manufacturing costs. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "Communications systems or networks generally comprise channels for sending and receiving data, audio and/or video information. A number of wired communications systems or networks, such as cable, Ethernet and Gigabit Ethernet systems, generally include stationary channels. A stationary channel is generally one in which the channel response and the noise statistics do not vary significantly or regularly (if at all) over time. Techniques are available for such systems and networks to reduce the adverse effects of random and/or systematic fluctuations in communications parameters. However, such techniques generally do not affect communications parameters that tend to fluctuate over time.\nCertain communications systems or networks, such as powerline systems or networks, include cyclostationary channels. A cyclostationary channel is one in which the channel response, the noise statistics and/or channel attenuation vary periodically. Also, powerline channels, which use conventional AC power lines for communications, experience regular fluctuations in noise and signal attenuation. The period of the noise power function in a powerline channel is generally the inverse of twice its frequency. In a typical case, that period is ( 1/120 Hz), or 8.3 msec. Noise in a powerline channel may also be introduced by other sources. During “spikes” in the noise power in a powerline channel, the signal-to-noise ratio (SNR) can be reduced sufficiently to cause errors in the data. Particularly when the signal strength is low, a low SNR can cause significant reliability problems.\nReferring now to FIG. 1, the powerline transmits a 60 Hz or 50 Hz sinusoidal power signal 10 supplying AC. The noise 20 in power lines may change within a period TAC or TAC/2, where TAC is the duration of one cycle of the AC power supply and/or voltage, typically 1/60 or 1/50 seconds (see, e.g., “Modeling of Cyclostationary and Frequency Dependent Power-Line Channels for Communications,” by Katayama et al.). As shown in FIG. 1, under common normal operating conditions, powerline channel noise 20 typically comprises a “zero power” component 30 (representative of the noise in the channel at the points where the AC power curve 10 crosses the zero power axis 15) and a burst noise component 40 (representative of the noise in the channel where the AC power curve 10 is at or near its maximum value[s]), although noise in cyclostationary channels may vary, depending upon operating conditions, equipment, power supply variations, etc. Thus, under normal operating conditions, cyclostationary channel noise 20 may be somewhat related to the absolute value of the AC power signal 10.\nA typical packet 100 for burst mode transmission is shown in FIG. 2. The packet 100 typically comprises a preamble 110, a synchronization sequence (or “syncmark”) 120, and data 130. The synchronization sequence 120 is used during frame synchronization. Under burst noise, especially cyclostationary burst noise, there is a relatively high probability that bits in the synchronization sequence 120 may be corrupted. For 10 Kb/s DBPSK modulation (0.1 ms/bit), if the cyclostationary noise has a burst length of half of its cycle (i.e., half of TAC/2=4.15 ms for 60 Hz AC power), about 42 bits could be corrupted by cyclostationary burst noise. This implies that the synchronization sequence 120 should be longer, and perhaps significantly longer, than 42 bits.\nIn systems without burst noise, the preamble 110 is often a 010101 sequence, and it is often used in carrier recovery. Due to the long, cyclostationary burst noise on the powerline, it is possible that the entire preamble 110 can be corrupted by burst noise, which means a very long preamble (e.g., more than 4.2 ms, or 42 bits at a 10 Kb/s rate) is desirable in a 60 Hz powerline channel to improve the success rate in carrier recovery. However, it is naturally desirable to minimize the number of bits dedicated to non-data (e.g., control and/or identification) functions, particularly over a sometimes noisy and/or problematic medium such as power lines. In some applications (e.g., voice transmissions), one might wish to completely avoid transmitting information during periods of high noise, and thus, one might keep the packet length small to fit packets into the “low noise” periods on the power line.\nA need therefore exists to improve communications in channels in which time-dependant fluctuations are a potential source of error or corruption, in order to reduce errors, reduce shutdowns, increase the success rate of such communications and/or increase communications uptime in networks that include such channels."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid droplet jetting apparatus which jets a liquid droplet and an ink-jet printer which jets an ink.\n2. Description of the Related Art\nAn ink-jet head which includes a plurality of pressure chambers communicating with a plurality of nozzles respectively, and a common liquid chamber supplying an ink to these pressure chambers, and which jets a droplet of the ink from a nozzle has hitherto been known. Such an ink-jet head selectively applies a pressure to the ink in the pressure chambers, and jets a droplet of the ink on to a recording paper from a nozzle communicating with the pressure chamber which is selected. For example, a serial ink-jet head described in Japanese Patent Application Laid-open No. H10-291311 includes a plurality of pressure chamber rows (three pressure chamber rows), and each pressure chamber row includes a plurality of pressure chambers (ink cavities) arranged in a longitudinal direction (paper feeding direction) of the serial ink-jet head. Furthermore, a common liquid chamber (ink supply chamber) which supplies the ink to each of the pressure chambers is extended along the longitudinal direction of the serial ink-jet head, in an area on an outer side of the pressure chamber rows and an area on an inner side (between the adjacent pressure chamber rows) of the pressure chamber rows.\nThe common liquid chamber which is arranged at the inner side of the two adjacent pressure chamber rows supplies ink to each of the pressure chambers belonging to the two pressure chamber rows positioned at both sides thereof. Therefore, in a case in which it is necessary to jet the ink simultaneously from a multiple number of pressure chambers, since an amount of ink supplied to the pressure chambers from the common liquid chamber at the inner side is increased temporarily, there is a possibility that the sufficient amount of ink is not supplied to each pressure chamber. Therefore, in the ink-jet head described in Japanese Patent Application Laid-open No. H10-291311, an ink communicating channel which communicates the common liquid chamber at the inner side with the common liquid chamber at the outer side, is provided in a partition wall (a column) which partitions the pressure chambers belonging to the each pressure chamber row. Moreover, when the amount of ink which is supplied from the common liquid chamber at the inner side, to the pressure chamber is increased temporarily, the ink is replenished from the common liquid chamber at the outer side to the common liquid chamber at the inner side via the ink communicating channel. Therefore, an insufficiency of the ink in the common liquid chamber at the inner side is prevented to a possible extent."}
{"text": "Electronic ovens heat items within a chamber by exposing them to strong electromagnetic fields. In the case of typical microwave ovens, the electromagnetic fields are a result of microwave radiation from a magnetron, and most often take the form of waves with a frequency of either 2.45 GHz or 915 MHz. The wavelength of these forms of radiation are 12 cm and 32.8 cm respectively. While heating, the electromagnetic waves in the chamber of a magnetron-powered microwave oven may drift or hop in frequency for short periods of time, generally within a range of +/−5%. For purposes of this disclosure, the mean temporal wavelength of an electromagnetic wave is referred to as the “dominant wavelength” of the associated electromagnetic wave, and dimensions of an electronic oven that are given with respect to a frequency or wavelength of an electromagnetic wave refer to the frequency or wavelength of the dominant wavelength of that electromagnetic wave.\nThe waves within the microwave oven that are not absorbed by the heated item reflect within the chamber and cause standing waves. Standing waves are caused by the constructive and destructive interference of waves that are coherent but traveling in different directions. The combined effect of the reflected waves is the creation of local regions of high and low microwave field intensity, or antinodes and nodes. The waves may interfere destructively at the nodes to create spots where little or no energy is available for heating. The waves interfere constructively at the antinodes to create spots where peak energy is available. The wavelength of the radiation is appreciable compared to the length scales over which heat diffuses within an item during the time that it is being heated. As a result, electronic ovens tend to heat food unevenly compared to traditional methods.\nElectronic ovens are also prone to heat food unevenly because of the mechanism by which they introduce heat to a specific volume of the item being heated. The electromagnetic waves in a microwave oven cause polarized molecules, such as water, to rotate back and forth, thereby delivering energy to the item in the form of kinetic energy. As such, water is heated quite effectively in a microwave, but items that do not include polarized molecules will not be as efficiently heated. This compounds the problem of uneven heating because different portions of a single item may be heated to high temperatures while other portions are not. For example, the interior of a jelly doughnut with its high sucrose and water content will get extremely hot while the exterior dough does not.\nTraditional methods for dealing with uneven cooking in electronic ovens include moving the item that is being heated on a rotating tray and homogenizing the distribution of electromagnetic energy with a rotating stirrer. These approaches prevent an antinode of the electromagnetic waves from being applied to a specific spot on the item which would thereby prevent uneven heating. However, both approaches are essentially random in their treatment of the relative location of an antinode and the item itself. They also do not address the issue of specific items being heated unevenly in the microwave. In these approaches, the heat applied to the chamber is not adjusted based on the location, or specific internal characteristics, of the item being heated."}
{"text": "1. Field of the Invention\nThe invention relates to a low pressure dew and frost point indicator and more particularly to the cooling module of such an indicator.\n2. Description of the Prior Art\nThe prior art known to the applicants is best represented by the prior Sheldon U.S. Pat. Nos. 3,874,220 and 3,886,784, Ford U.S. Pat. Nos. 3,374,658, 3,152,475 and 3,460,373 and Hankison U.S. Pat. No. 3,589,165.\nIn prior devices, there has been no positive and simple way of positioning the frost point indicating mirror correctly with respect to the cold air stream which is directed against it.\nPrior mirrors have also been subject to temperature gradients and undesirable changes of temperature caused by contact with adjacent module components and/or subjection to atmospheric or high pressure gases entering from the surrounding area.\nEjectors of the type used in the prior art have not had the capability of being adjusted for optimum operation of the cold air generation system and then being pre-set in optimum setting.\nIn case a new reading is required, necessitating raising the temperature in the sampling chamber, the prior art has utilized purge or blow down valves to clear the sample chamber or has cut off air controlled by the input to the sample chamber with relatively large valves in order to allow the sample chamber to become warmer."}
{"text": "1. Field of the Invention\nThis invention relates to a novel surgical stapling instrument, particularly useful in intralumenal anastomosis and characterized in affording semi-disposable properties. This invention particularly relates to a surgical stapling instrument employed in placing two or more evenly spaced, laterally aligned rows of staples in an internal organ during a surgical procedure.\nSuch instruments generally comprise two elongate jaw members, respectively affording a staple holder and an aligned anvil for closing the staples. The staple holder is designed to afford at least two rows of spaced staples on either side of a longitudinal opening in the anvil member. A slider assembly including a knife is provided to move longitudinally along the jaw members to eject and close staples while severing the tissue by passage of the knife through the longitudinal opening in the anvil member.\nIt is noted at the outset that the prior art sometimes uses the term \"pusher bar\" to describe what is herein referred to as the \"slider assembly.\" We further note that the term \"pusher bar\" is herein used to describe a component of the disclosed invention that is distinctive from the \"slider assembly.\"\n2. Description of Related Art\nIn earlier prior art instruments, the staple holders, associated anvil members, and slider and knife assemblies have typically been constructed of disposable plastic materials and low-cost metals. Frame members of the instruments have been made of more durable materials suitable for sterilizing so that the instruments could be available for repeated use. By way of example, one instrument of this type is disclosed in U.S. Pat. No. 3,499,591, the disclosure of which is incorporated herein by reference.\nMore recent prior art instruments have been designed to be wholly disposable. This eliminates the need for repeated sterilization and insures that the instrument is ready for use without concern for necessary preparatory steps. Inasmuch as the instrument is to be discarded after one surgical operation, its design typically employs readily available, low-cost materials of construction to achieve an economically attractive result.\nIn such stapling instruments, large forces are generated when the tissue is clamped between the jaws, the staples are made to pass through the tissue, and the staples are closed in contact with aligned depressions in the anvil. These forces are of such a magnitude as to typically cause vertical separation and lateral distortion of the jaws and thus affect correct staple closure. Such forces become a primary concern when lighter materials are used, as, for example, in the design of the disposable instrument.\nEarlier prior art instruments were designed to provide that the knife blade was an integral part of the sterilizable instrument so that the knife blade was resused when the unit was reloaded with staples for future use. With time, the knife blade became dulled by the repeated stapling procedures so that cuts were no longer as clean as with an unused blade.\nMore recent prior art, providing fully disposable surgical staplers, as suggested by U.S. Pat. No. 4,429,695, does not provide for replacement of the staple holder. Accordingly, the instrument must be abandoned after one firing process even though the operative surgical procedure may call for more staples or for more than one stapling operation. The need for a second instrument in a single patient's operation leads to more expensive procedures.\nFurther, support shoes which are carried by the slider and knife assemblies to provide vertical and lateral support to the jaws regularly encounter friction forces that make pushing difficult and do not permit the smooth, sliding action helpful to successful stapling procedure.\nRecent prior art instrument designs generally require the use of a knife blade because the support shoes are connected to a plate member which passes through the tissue and thus must follow behind the knife blade. The support shoes in these prior art instruments restrain the jaws from vertical separation and lateral distortion.\nIn the prior art, the slider assembly (sometimes referred to as the \"pusher bar\" in the prior art) is designed to provide a thumb tab projecting laterally outward from between the upper and lower frame members on only one side of the instrument so that the propulsion of the pusher bar is effected by one hand laterally displaced from the longitudinal movement of the pusher bar. This action may place disorienting strains on the normal movement of the slider assembly and the instrument and may also affect the placement of the staples in the tissue.\nThere remains a need for an improved surgical stapler instrument, employing economic materials of construction, which can provide for the affording of sufficient staples to complete a patient's surgical operation, for the minimization of forces tending to hinder proper operation of the instrument, and for more convenient manipulation during the surgical operation procedure. Additionally, there remains a need for making the inclusion of a knife unit optional, as a matter of economy where one is not needed."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display (LCD), and more particularly, to a response time accelerator system and method for driving a liquid crystal display.\n2. Description of the Related Art\nCurrent LCD technologies have problems associated with liquid crystal response time. Thus, because the response time of the liquid crystal forming pixels in an LCD panel is relatively slow, a user sees an after image, when a TV displays a large number of moving images.\nEach liquid crystal cell within an LCD panel allows light to pass through or blocks the light when an applied bias voltage rotates a lattice. Usually, liquid crystals cannot respond in real-time to data changes since the time required to respond to the bias voltage is on the order of tens of milliseconds and bias voltage applied across the crystal varies from frame to frame (1/75 seconds at SXGA resolution). For example, if the bias voltage applied across the liquid crystal within an LCD panel is set to a gray level of 255 for 8-bit image data, the brightness level after the actual response of the liquid crystal is less than 255, which causes vertical stripe patterns to generate an after image. The liquid crystal response characteristics is aggravated as the frame period decreases due to increased resolution.\nThe most commonly used method for preventing liquid crystal panel after image retention due to a slow response is to appropriately preprocess image data prior to processing it in a source driver for driving a liquid crystal panel. To implement this technique, a response time accelerator (RTA) has been adopted. However, the RTA has many problems.\nIn order for the liquid crystal to keep up with current frame data, the RTA basically compares a previous frame data with the current frame data, interpolates the two frame data and finds a new current frame data with respect to which a response time can be accelerated according to the result of the comparison. The current data input to the RTA in real time is compared with the previous frame data stored in a memory such as SDRAM outside the RTA.\nA method for accelerating a response time applied to a conventional RTA involves storing four to six most significant bits (MSB's) of RGB (Red, Green, Blue) data as frame data in an external memory such as SDRAM. However, when storing only the MSBs (e.g., n MSBs) in an external memory such as SDRAM as frame data, a pixel data truncation error (PDTE) occurs. If only the MSBs are used as previous frame data Pn−1, the number of gray level levels that can be compared between the previous frame data Pn−1 and current frame data Pn is (2n×2n) where n bits Pn−1[7:8−n] of previous frame data Pn−1 and n bits Pn[7:8−n] of current frame data Pn are used. This causes an error to occur. That is, a quantization error occurs because the number of gray levels available for comparison is reduced to (2n×2n) from “256×256”, which is the number of gray levels available when 8-bit data is entirely used for each frame data. Here, the term “quantization error” is given because the error occurs sporadically due to unused 8-n LSBs.\nFor example, assuming a value derived by shifting n bits Pn[7:8-n] of the current frame data Pn to the right by 8-n bit positions is K (e.g., Pn[3:0] if n=4), the quantization error occurs when outputting an overshoot value exceeding or an undershoot value below the current frame data to the panel at gray level K*2(8−N). Eventually, since the quantization error occurs periodically at gray level produced by multiplying 2(8−N) by the value derived by shifting current frame data for each gray level by 8-n bit positions, noise is generated at regular intervals in vertical stripes to be displayed on a liquid crystal panel. The above description, as previously mentioned, is under the assumption that the LCD uses 8-bit RGB data for each pixel making it possible to display 256 gray levels.\nHere, the new current frame data found by interpolation of the previous frame data Pn−1 and the current frame data Pn is usually stored in a table memory within an RTA. The tabular values stored in the built-in table memory for interpolation, defining overshoot exceeding or undershoot below the current frame data, are experimentally determined based on the liquid crystal characteristics of a panel. For example, given the fact that the response time of liquid crystal in the panel is relatively slow, if the current frame data is greater than the previous frame data, the new current frame data is assigned a tabular value greater than actual data, which is then output to the panel. Similarly, if the current frame data is less than the previous frame data, the current frame data is assigned a tabular value less than the actual data, which is then output to the panel.\nAnother drawback of the method for accelerating a response time implemented using a conventional RTA is that it limits the amount of overshoot and undershoot for each output data to a gray level range from 0 to 255 since the number of bits of RGB image data is limited to 8. Therefore, if the current frame data Pn has an extremely large value (around a gray level of 255) or an extremely small value (around a gray level of 0), the output data does not have sufficient amount of overshoot or undershoot. For example, if the current frame data Pn has a maximum gray level of 255, it is impossible to give an overshoot greater than the current frame data Pn to the output data since the overshoot is limited to the maximum gray level of 255. Thus, if data output from the panel is at a gray level of 255, the liquid crystal panel will return a value less than 255 in response. Such an overshoot or undershoot limitation makes it difficult to improve a response time."}
{"text": "The UWB-RF channel has recently become available in the USA (e.g., in the frequency range 3.1 GHz-10.6 GHz allocated for indoor as well as outdoor use). European and Japanese authorities are preparing similar rules to enable commercial marketing and use of UWB-RF devices.\nShort-range wireless technologies in the wireless local area network (WLAN) space as well as the wireless personal and body area network spaces (WPAN/WBAN) continue to proliferate rapidly. In particular, wireless links based on infrared (IR or Ir) light emission have experienced an enormous growth in the recent past owing to ease of use and low cost. The Infrared Data Association (IrDA) has projected a total shipment of 1.3 Billion units by the year 2003, with over 250 Million IrDA devices installed in mobile phones and personal digital assistants (PDAs). This high adoption rate reflects the presence of a large number of IrDA applications in the marketplace and a sizable investment by the industry and users alike. On the other hand, radio-based short-range wireless devices such as Bluetooth™ (Bluetooth is a trademark owned by Bluetooth SIG, Inc.) are also deployed at a rapid pace, thereby reducing the relative size of the IrDA-based market segment rather quickly. This will influence the IR-based industry which is looking for alternate ways to leverage their investments and continue their business.\nTwo dominant technologies in the short-range wireless space—Bluetooth and IrDA —were specifically designed for replacing short-range wired communication links. Bluetooth is based on an RF technology offering less than 1 Mb/s of user data rate and IrDA offers a number of data rates, e.g., SIR (115 Kb/s), FIR (4 Mb/s), and VFIR (16 Mb/s). IrDA uses optical-IR transmission which requires line of sight (LOS) between transmitter and receiver. IrDA and Bluetooth technologies are also used for consumer-based applications within the single office/home office (SOHO) environment. Each has its own strengths and weaknesses with respect to link range, data rates, required regulation and costs.\nWhile some of these technologies' intended uses and applications overlap, each of the two systems offer all of the hardware and software layers that constitute a set of communication protocols. IrDA is an established and well proven point-to-point, narrow-angle, data-transmission standard that operates at bit rates between 9600 b/s and 16 Mb/s over distances of 0-1 m (mostly 20 cm). It has a wide range of supported hardware and software platforms and is well introduced in the market place. Though IrDA's benefits seem to abound, the technology has its disadvantages. For many applications, its limited range and pointing angle is rather inconvenient and its optical signals fail to penetrate walls or propagate around obstructions (IrDA was not designed to rely on diffused optical signal propagation). Additionally, IrDA devices must be nearly stationary to achieve synchronization with other IrDA devices (narrow optical beam). Where IrDA fails to be user friendly, radio-based systems offer advantages, such as non-line-of-sight (NLOS) transmission through walls and other obstructions. However, many RF-based systems, such as Bluetooth, are regulated differently in various countries worldwide and thus suffer non-trivial insertion into the markets.\nThe two main markets for short-range wireless devices are the portable and desktop computing and handheld markets. Due to architectural and range constraints, the level of adoption of IrDA in the desktop space has been relatively low. Current desktop designs in the marketplace are more appropriate for placing the desktop in a cabinet or on the floor. However, such a usage scenario is not favorable for point-and-shoot optical link applications as offered by IrDA.\nThe disadvantages of optical-IR communication links and the investments from the manufacturers and the user community in IrDA-based technology and applications could both be leveraged and maintained into the future by replacing or enhancing the optical transceiver of the IrDA system (i.e., part of the PHY layer) with a suitable radio-based transceiver.\nGiven the recent introduction of a license-free spectrum band between 3.1 GHz and 10.6 GHz (USA/FCC) for UWB radio devices, a new and inherently compatible radio technology has become available to design and implement such radio-based systems. As a result, the many applications and protocol stacks, including medium access control protocols, originally conceived for wireless optical links should be useable on the UWB radio channel.\nIt would be advantageous if the two types of communication systems could be brought together to make use of the wireless optical system and the radio channel. Owing to the different transmission media used (optical and radio frequency, respectively) as well as because of legacy considerations, these different wireless communication systems have traditionally also been designed with different medium access control (MAC) mechanisms; in addition, their positioning in the market place have been different in general, although their applications do overlap in certain areas (e.g., PDA data base synchronization with PC data base)."}
{"text": "The present disclosure relates generally to product packages, and relates more particularly to packages that include a tamper-evidence feature.\nWhen they purchase a product from a retailer, consumers like to know that the package containing the product has not previously been opened. This is particularly true with certain types of products such as foods or other products used in or on one's body. Accordingly, it is common for product packaging to include some kind of tamper-evidence device that is supposed to indicate in a readily visible manner whether or not the package has been opened. Various types of tamper-evidence devices have been developed, such as shrink bands, pressure sensitive adhesive tabs that tear through upon opening, laser scored areas that propagate a tear in the packaging material, and devices based on delamination of inks and coextruded films.\nThere remains a need for a tamper-evident package having the features of the package disclosed herein."}
{"text": "This invention relates generally to the formation of circuit boards or cards or the like, and more particularly to the formation of circuit boards or cards having 2 signal planes and one power plane (2S/1P) wherein the power plane is sandwiched between two layers of photopatternable dielectric material and on which layers of circuitry for the signal planes is disposed."}
{"text": "1. Field of the Invention\nThis invention relates to a grommet, and more particularly relates to a grommet that is mounted on a wire harness to be arranged in a motor vehicle and is attached to a burred through-hole in a vehicle body panel.\n2. Description of the Related Art\nA wire harness that extends from an outer compartment of a motor vehicle, such as an engine compartment, to an inner compartment, such as a driver's compartment, passes through a through hole in a vehicle body panel that partitions the outer and inner compartments. A grommet is mounted on the wire harness and is attached to the through-hole to effect waterproof, dustproof, and soundproof functions.\nThe prior art includes various kinds of grommets, including, a one-body grommet made of a rubber or an elastomer, and a two-body grommet that has a grommet main body made of a rubber or an elastomer and a resin inner incorporated with the grommet main body.\nJP HEI 9 (1997)-27226 A and FIGS. 6A to 6B herein disclose a two-body grommet that is used when a through-hole in a vehicle body panel is a burred hole. With reference to FIG. 6A, a two-body grommet 100 includes a grommet main body 110 made of a rubber and a resin inner support 120. The grommet main body 110 includes a waterproof sealing portion 111 that contacts a vehicle body panel P. On the other hand, a flange 122 projects from an end of an inner side cylindrical portion 121 of the resin inner support 120 and is inserted into an annular recess 112 in the grommet main body 110. A latch pawl 123 is formed by turning the other end of the inner side cylindrical portion 121 outward. The latch pawl 123 contacts and latches with a burr B projecting from a peripheral edge around a through-hole H in the vehicle body panel P.\nAn O-ring-like portion 111b is provided on the distal end of the waterproof sealing portion 111 of the grommet main body 110 and has a ball shape in cross section. The ball-like ring portion 111b is pressed onto the vehicle body panel P to effect a waterproof action. JP HEI 9 (1997)-27226 also discloses an example with another O-ring-like waterproof sealing portion 115 inside the waterproof sealing portion 111 (see FIG. 6B) and an example in which a lip-like waterproof sealing portion 116 is provided outside the O-ring-like waterproof sealing portion 115 (see FIG. 6C). These two waterproof sealing portions 115 and 116 project from an outer peripheral wall at the same height from the annular recess 112 and are oblique to the wall.\nAs shown in FIGS. 6B and 6C, the conventional grommet includes two waterproof sealing portions with equal projecting heights H1. One waterproof sealing portion can exert its waterproof function even if the other waterproof sealing portion bites foreign particles and curls up to lose its waterproof function.\nAs shown in FIG. 7A, a burred hole formed in the vehicle body panel P has a burr B that projects from an inner peripheral edge around the through-hole H in the vehicle body panel P. A height (a) of the burr B is likely to be random, as shown in FIGS. 7B, 7C, and 7D. On the other hand, a height H3 between the latch pawl 123 of the resin inner support 120 and a lip 111r on a distal end of the waterproof sealing portion 111 of the grommet main body 110 are constant.\nAccordingly, as shown in FIG. 7C, if the projecting height (a) of the burr B is formed into a designed height precisely, the lip 111r of the waterproof sealing portion 111 is pressed onto the vehicle body panel P by a given contact pressure. However, as shown in FIG. 7B, if the height (a) of the burr B is small, the lip 111r of the waterproof portion 111 lightly contacts the vehicle body panel P and is not pressed onto the panel P by the given contact pressure. This may cause a clearance between the lip 111r and the panel P. As shown in FIG. 7D, if the projecting height (a) of the burr B is great, a contact wall 111a of the waterproof sealing portion 111 contacts the vehicle body panel P and the lip 111r does not push the panel P. This will cause a clearance between the lip 111r and the panel P or this will exert less than the given contact pressure between the lip 111r and the panel P.\nAccordingly, even if two waterproof sealing portions are provided on the grommet main body, as shown in FIGS. 6B and 6C, they cannot reply to the random projecting heights of the burrs, since the projecting heights of the two waterproof sealing portions are same.\nIn view of the above problems, an object of the invention is to provide a grommet that can provide a high waterproof function at a burred through-hole in a vehicle body panel even if projecting heights of the burrs are random."}
{"text": "The exposure of phosphatidylserine (PS) and other aminophospholipids (aminoPL) on the surface of activated or injured blood cells and endothelium is thought to play a key role in the initiation and regulation of blood coagulation. De novo surface exposure of aminophospholipids has also been implicated in the activation of both complement and coagulation systems after tissue injury, and in removal of injured or apoptotic cells by the reticuloendothelial system. Although migration of these phospholipids (PL) from inner-to-outer plasma membrane leaflets is known to be triggered by elevated intracellular [Ca2+] ([Ca2+]) and to be associated with vesicular blebbing of the cell surface, little is known about the cellular constituents that participate in this process.\nAs described in Ser. No. 08/790,186, cell surface PS has a role in coagulation, programmed cell death and clearance by the reticuloendothelial system. Ser. No. 08/790,186 also describes regulation of the transmembrane distribution of PS, the role of calcium in the collapse of phospholipid asymmetry, and the role PL translocation in Scott Syndrome.\nBassxc3xa9, et al. and Stout, et al. recently reported the purification and preliminary characterization of an integral RBC membrane protein that, when reconstituted in liposomes, mediates a Ca2+-dependent transbilayer movement of PL mimicking plasma membrane PL reorganization evoked upon elevation of [Ca2+]c (F. Bassxc3xa9, et al., J. Biol. Chem. 271:17205-17210, 1996; J. G. Stout, et al., J. Clin. Invest. 99:2232-2238, 1997). Evidence that a protein of similar function must also be present in platelets was recently reported by Comfurius, et al. (P. Comfurius, et al., Biochemistry 35:7631-7634, 1996).\nNeeded in the art is a method for modulating the activity of phospholipid scramblase within a cell, organ or tissue in which one wishes either to reduce the potential for thrombosis, clot formation, or cell clearance (by decreasing cellular PL scramblase expression or activity) or to promote hemostasis or cell clearance (by increasing cellular PL scramblase expression or activity).\nThe present invention relates to the creation and use of antithrombotic and thrombostatic reagents that rely on the properties of a protein preparation that mediates Ca2+-dependent transbilayer movement of membrane phospholipids.\nIn one embodiment, the present invention is a preparation of a plasma membrane phospholipid scramblase (xe2x80x9cPL scramblasexe2x80x9d). Preferably, the protein is approximately 35-37 kD as measured on a 12.5% SDS-polyacrylamide gel under reducing conditions. In a most preferred form of this invention, the preparation is a human or a mouse PL scramblase.\nIn one preferred embodiment of the present invention, the PL scramblase comprises SEQ ID NO:2, representing human PL scramblase, possibly with conservative or functionally equivalent substitutions.\nIn the most preferred embodiment of the present invention, the PL scramblase, preferably comprising SEQ ID NO:2, has been modified by the action of mammalian cellular enzymes to covalently incorporate phosphorous at one or more Thr, Ser, or Tyr residues or a fatty acid, preferably palmitate, at a cysteine residue.\nThe present invention is also a preparation of a murine cell protein, wherein the protein is a plasma membrane phospholipid scramblase, preferably wherein the protein is approximately 35 kD as measured on a 12.5% SDS-polyacrylamide gel under reducing conditions.\nIn one preferred embodiment of the present invention, the murine PL scramblase comprises SEQ ID NO:4, possibly with conservative or functionally equivalent substitutions.\nIn the most preferred embodiment of the present invention, the murine PL scramblase comprising SEQ ID NO:4 has been modified by the action of mammalian cellular enzymes to covalently incorporate phosphorous one or more Thr, Ser, or Tyr residues, and a fatty acid, preferably palmitate, at a cysteine residue.\nThe present invention is also a DNA sequence encoding the PL scramblase. Preferably, this DNA sequence comprises SEQ ID NO:1. Most preferably, this DNA sequence comprises residues 211-1164 of SEQ ID NO:1.\nThe present invention is also a DNA sequence encoding the murine PL scramblase. Preferably, this DNA sequence comprises SEQ ID NO:3. Most preferably, the DNA sequence comprises residues 192-1112 of SEQ ID NO:3.\nIn another embodiment, the present invention is a method of preventing the surface exposure of plasma membrane phospholipids and reducing the procoagulant properties of a cell by delivering to the cell a mutant phospholipid scramblase. This scramblase is preferably mutated at a site of post-translational modification, most preferably the site is selected from the group consisting of Asp273-Asp284, Thr161 and Cys297 of human PL scramblase SEQ ID NO:2 or the corresponding conserved residues in mouse or other mammalian PL scramblase.\nIn one embodiment, a gene construct encoding a mutant phospholipid scramblase is delivered to the cell. In an alternative embodiment, the mutant protein itself is delivered.\nThe present invention is also an inhibitor of the PL scramblase activity of PL scramblase. This inhibitor may be an antisense nucleotide derived from the DNA sequence of PL scramblase. In another embodiment, the inhibitor is a peptide sequence that is a competitive inhibitor of PL scramblase activity. In another embodiment, the inhibitor is an antibody, preferably a monoclonal antibody, raised against PL scramblase.\nIn another embodiment, the inhibitor works by modifying or inhibiting the post-translational modifications of the PL scramblase that are disclosed below in the Examples. For example, analysis of the primary PL scramblase sequence reveals a potential site of phosphorylation by protein kinase C or other cellular kinase (Thr161), a potential site for acylation by fatty acid (Cys297), and a potential binding site for Ca2+ ion provided by an EF-hand-like loop spanning residues Asp273-Asp284. These residues and motifs are conserved 25 in the mouse PL scramblase.\nIn one embodiment, the inhibitor is a compound that prevents thioacylation of the protein.\nIn another embodiment, a mutant phospholipid scramblase is provided in which cysteine residues, preferably Cys297 of SEQ ID NO:2 (or the equivalent residue in the conserved region of another PL scramblase), have been replaced by alanine or other non-conservative substitution.\nThe present invention is also a method for preventing the surface exposure of plasma membrane phosphatidylserine, phosphatidylethanolamine and cardiolipin on the surface of in vitro stored leukocytes, lymphocytes, platelets or red blood cells. This method comprises the steps of adding an inhibitor of PL scramblase activity to the stored blood cells.\nThe present invention is also a method for prolonging survival of transplanted organs comprising the step of adding an inhibitor of PL scramblase activity to an organ perfusate during in vitro organ storage. The present invention is also a method for prolonging the survival of transplanted cells, tissues, and organs by genetically engineering the cells to be transplanted so as to alter their expression of plasma membrane PL scramblase in order to reduce exposure of PS and other thrombogenic phospholipids at the plasma membrane surface, thereby reducing the risk of infarction due to fibrin clot formation.\nThe present invention is also a method for prolonging the in vivo survival of circulating blood cells (erythrocyte, platelets, lymphocyte, PMN\"\"s, and monocytes) comprising the step of preventing surface exposure of plasma membrane phosphatidylserine on the surface of the cells by exposing the blood cells to an inhibitor of PL scramblase activity.\nThe present invention is also a method for preventing the procoagulant activities of erythrocytes in sickle cell disease comprising the step of inhibiting erythrocyte PL scramblase in a sickle cell patient.\nThe present invention is also a method for treating autoimmune and inflammatory diseases comprising the step of treating a patient with an inhibitor of the PL scramblase activity of PL scramblase.\nThe present invention is also a method for diagnosing individuals with reduced or elevated capacity for platelet-promoted or erythrocyte-promoted fibrin clot activity comprising the step of quantifying the cellular expression of PL scramblase. This quantitation may take the form of immunoblotting using an antibody to PL scramblase, an ELISA assay using an antibody to PL scramblase, flow cytometric analysis of the binding of monoclonal antibody reactive against the predicted extracellular domain of PL scramblase (residues Ser310-Tryp318 of sequence disclosed in SEQ ID NO:2 or the equivalent residue in the conserved region of another PL scramblase) or using oligonucleotides derived from PL scramblase cDNA and the polymerase chain reaction. In one method of the present invention, the quantitation is performed by isolating PL scramblase from a patient blood sample, measuring the amount of PL scramblase isolated and comparing the measurement with a control sample. The measurement may be by isolating PL scramblase from a patient blood sample and measuring via densitometry the amount of PL scramblase protein electrophoresed in a stained electrophoretic gel.\nIt is an object of the present invention to provide a preparation of a PL scramblase.\nIt is another object of the present invention to genetically alter the level of expression of PL scramblase by delivery of cDNA representing sense or antisense nucleotide sequence ligated into a suitable mammalian expression vector.\nIt is another object of the present invention to provide an inhibitor of PL scramblase PL scramblase activity.\nIt is another object of the present invention to provide an antithrombotic agent.\nIt is another object of the present invention to create cells, tissue, and organs for transplantation that have increased potential for survival and reduced potential for causing fibrin clot formation and vascular thrombosis when grafted into a recipient host.\nIt is another object of the present invention to create an animal, preferably a mouse or pig, that has been genetically engineered so that the PL scramblase gene is not expressed.\nOther objects, advantages and features of the present invention will become apparent after one of skill in the art reviews the specification, claims and drawings herein."}
{"text": "This invention relates to a valve mechanism for an internal combustion engine and, more particularly, to an improved cylinder head and valve actuating system for an overhead camshaft internal combustion engine.\nIt is well recognized that the performance of internal combustion engines can be improved through the use of overhead valves operated by overhead mounted camshafts. When an overhead camshaft arrangement is employed, the cylinder head assembly trends to become much more complicated than with a conventional pushrod operated engine. This is because the cylinder head must, in addition to supporting the valves and the valve springs, provide support for the camshaft and for the actuators for the valves. Frequently, the valves are directly operated and this means that an arrangement must be provided for slidably supporting the thimble tappets that operate the individual valves from the camshaft.\nIf all of these functions are performed primarily by the cylinder head, then the cylinder head casting becomes extremely complicated. In addition, since it is necessary to machine the bearings for the camshaft provided by the cylinder head, and the bores for slidably supporting the tappets, than machining operations also add significantly to the cost of the cylinder head. Furthermore, there is the problem of assembly of all of components into such a unitary cylinder head assembly, and the problems of accessibility the various fasteners for securing the cylinder head to the cylinder block and the bearing caps to the cylinder head for journaling the camshaft. Of course, the problems mentioned above are complicated when the engine employs twin overhead camshafts.\nIt has been proposed, therefore, to employ a construction wherein the main cylinder head member itself does not have to perform all of these functions. For example, in U.S. Pat. No. 4,612,885, entitle Camshaft Bearing Arrangement For Overhead Cam Engine, issued Sep. 23, 1986, in the name of Masaaki Yoshikawa, and assigned to the assignee hereof, there is depicted a cylinder head arrangement wherein the main cylinder head assembly only supports the poppet valves for their movement and the return springs for the poppet valves. The camshafts and valve actuating tappets are supported in a separate cam carrier that is affixed to the cylinder head and thus can be machined and cast separately simplifying the aforenoted problems. However, with the arrangement shown in that Patent, the cam carrier forms the outer periphery of the cylinder head and the cam cover must sealingly engage it. In addition, the cam carrier must have a sealing arrangement around its outer periphery with the upper surface of the cylinder head to afford sealing. Hence, substantial addition machining operations are required.\nIt is, therefore, a principal object of this invention to provide an improved cylinder head assembly for an overhead camshaft internal combustion engine.\nIt is a further object of this invention to provide an improved cylinder head assembly for an overhead cam internal combustion engine wherein at least some of the valve actuating tappets and at least one of the camshafts are supported by a separate cam carrier member that is affixed to the cylinder head but in such a way that this cam carrier member need perform no sealing functions for the overall cylinder head assembly.\nIt is a further object of this invention to provide an improved and simplified cylinder head assembly for an internal combustion engine having an overhead camshaft wherein machining and assembly operations are considerably simplified.\nIn addition to the problems already noted, it is also desirable to ensure that the cam actuating tappets and camshafts are well lubricated. When a separate cam carrier member is provided, this can present additional difficulties in insuring that the camshafts are adequately lubricated as are the valve actuating tappets.\nIt is, therefore, a still further object of this invention to provide an improved cylinder head assembly including an improved arrangement for lubricating camshaft journals thereof."}
{"text": "Polyaspartic acid is a peptide chain in which amide linkages extend the chain. In the thermal polymerization of aspartic acid, the stereochemistry of the aspartic acid is racemized and the formation of both .alpha. and .beta. carboxylic acid groups have the ability to react to form such amide bonds. Such materials have been used for fertilizers and scale inhibition agents. They are particularly useful for the prevention of scale deposition in boiler water, reverse osmosis membranes, detergents and as inhibitors of dental tartar and plaque formation (tartar barrier agents). These materials are ready biodegradable. Methods for the preparation of polyaspartic acid have been developed (See U.S. Pat. Nos. 5,057,597 and 4,839,461 and U.S. patent application Ser. No. 07/882,919, filed May 14, 1992, Louis L. Wood, and U.S. patent application Ser. No. 07/926,242, filed Aug. 7, 1992, Louis L. Wood).\nBiodegradability, calcium ion exchange ability and the disruption of calcium salt crystal structure are important properties of materials used in the prevention of scale deposition in boiler water, on reverse osmosis membranes, in detergent use and as inhibitors of dental tartar and plaque formation (tartar barrier agents). We searched for economically useful materials, having a greater retention on the object wherein inhibition of scale deposition is desired. Other desirable properties were greater stability to biodegradation in addition to intrinsic value for the prevention of scale deposition in boiler water, on reverse osmosis membranes, during detergent use and as inhibitors of dental tartar and plaque formation (tartar barrier agents). We have found that the addition of polycarboxylic acids in the thermal polymerization of maleic acid or aspartic acid produced novel and highly effective copolymers which possessed these properties."}
{"text": "Swimming boards and swimming aids are known, e.g. from DE-PS 593 763, DE 43 33 941 A1, DE-GM 1 770 532 as well as the prospectus of the company, Volley Sportartikel GmbH, Ispo 1989, pages 18 to 20. DE-PS 593 763 discloses an auxiliary device for teaching oneself to swim, wherein a shaped board comprises, at the rear end, a semicircular recess for supporting the chin of the swimmer as well as, on its top surface, at least two handles which are provided at arm's length distance in forward direction on the top surface of the board.\nDE-GM 17 70 532 discloses a swimming board, which comprises a substantial body having at least two handles and a cutout for the head or chin of the swimmer. The swimming board is used as a simple aid for learning to swim, for lifesaving as well as an aid for rescuing shipwrecked persons.\nDE 43 33 341 A1 discloses a swimming aid comprising buoyancy elements, which in numerous variants are connected to a carrier element. The swimming aid is fastened by a belt to the waist region of the swimmer.\nSaid previously known swimming boards, aids for learning to swim and rescue aids are usable only for their, in each case, very narrowly defined intended application.\nSwimming systems are likewise known, particularly in the form of a two-part training aid for swimmers, which comprises a flat basic element and a resistance element insertable vertically into the latter.\nThe flat basic element is substantially rectangular with a trapezoidal moulding provided in the top region. The corners are in each case rounded-off. In the upper trapezoidal projection, an elongate rectangular slot is provided. In the lower region of the flat basic element, round openings are provided at a distance from the outside edge of the basic element. On the one hand, feet may be inserted through said round openings. On the other hand, however, the front portions of the hand may engage therein, with the balls of the thumbs being held firmly against the outside edge of the flat basic element. The flat resistance element has a middle region, which fits exactly into the elongate rectangular slot of the flat basic element. Outside of said middle region, the resistance element protrudes at the one side. In so doing, it forms a substantially rectangular region with rounded corners and two round openings arranged symmetrically relative to one another in the surface of said region. The other sub-region of the resistance element adjacent to the region engaging into the flat slot is provided in an inward direction with recesses, the dimensions of which are such that there is room inside for the ankles. The recess is delimited at the outer side of the resistance element by elongate end regions projecting in a hammer-like manner and describing a curve with a flat edge. The resistance element when inserted into the slot inside the basic element creates a water resistance, wherein said resistance element on account of the water flowing against it may execute a tilting motion inside the rectangular slot.\nSuch swimming systems, while being already somewhat more versatile than the swimming aids initially discussed at the start, are however still usable only in a very restricted and limited manner.\nThe object of the invention is to provide a swimming and/or gymnastics system according to the preamble of claim 1, which may be grasped and held fast very well by the hands, involves only a low expenditure of material and is usable in an extremely versatile manner for swimming exercises and gymnastics exercises in water but also on land."}
{"text": "1. Field of the Invention\nThis invention relates to a method for extracting the LCR (inductance, capacitance, and resistance) values of connecting line(s) in semiconductor integrated circuit (hereinafter LSI) design processes and to an LCR extraction program for causing a computer to execute that method, and particularly to a method and program capable of extracting highly precise LCR values of connecting line(s), taking into consideration the fluctuations in pattern width associated with the finer miniaturization of layout patterns.\n2. Description of the Related Art\nLSI design processes are ordinarily performed on computer CAD. An LSI design process has a logic design process for connecting logic gates to design a logic circuit, a layout design process for laying out that logic circuit on an actual chip, a process for extracting LCR (inductance, capacitance, and resistance) values of laid out connecting lines and determining the delay time for signal paths from the extracted LCR and the AC characteristics of cells and macros, a logic simulation process for checking whether the logic circuit will operate normally using that delay time, and a physical verification process for checking that the layout data satisfies the design rule.\nBased on t he layout design, layout data, including wiring pattern data for each layer on the chip, are produced. Based on these layout data, the connecting line LCR values are extracted. The LCR extraction process, delay time computation process, and logic simulation process are generally provided in one program module.\nIn the LCR extraction process described above, the connecting line resistance R, capacitance C, and inductance L are extracted, either by arithmetic computation or by referencing a parameter table, according to the width of the wiring, the distances to adjacent wiring, and the overlapping area contained in the layout data.\nIn conjunction with the finer miniaturization of patterns in recent years and due to the influence of manufacturing processes, the actually formed pattern width sometimes differs from the pattern width in the layout data. One manufacturing process thought to affect the pattern width is the connecting line etching process. In processes such as reactive ion etching (RIE), wherein a reaction gas is introduced in a high vacuum atmosphere, a plasma state is induced by applying a high frequency voltage, and a wiring layer of aluminum or the like is etched, a phenomenon is observed whereby the pattern width fluctuates, becoming thinner or thicker, in conjunction with finer pattern miniaturization.\nIn conjunction with this pattern width fluctuation phenomenon, with a conventional method for extracting LCR values from the layout data, problems are encountered in that accurate LCR values cannot be extracted, the precision of the signal propagation delay time on signal paths determined by the extracted LCR values is low, and suitable logic simulation cannot be effected.\nThereupon, an object of the present invention is to provide both a new method wherewith it is possible to extract LCR values of connecting lines and the like, taking pattern width fluctuation into consideration, and a computer program for implementing that method.\nAnother object of the present invention is to provide both a method wherewith accurate LCR value extraction can be performed, even when the LSI design rule has been made even finer, and a computer program for implementing that method.\nIn order to attain the objects noted above, in one aspect of the present invention, an LCR extraction method for extracting LCR values containing at least one of resistance, capacitance, and inductance, from layout data at least having wiring pattern data in a plurality of wiring layers, has the steps of:\ngenerating the LCR values, for a wiring pattern being looked at, based on the layout data for the wiring pattern being looked at;\nfinding the pattern congestion level in an area where the wiring pattern being looked at exists, based on the layout data; and\ncorrecting the LCR values when the pattern congestion level is higher (more congested) than a prescribed reference value, based on pattern fluctuation values depending on the pattern interval between the wiring pattern being looked at and an adjacent pattern.\nPattern width fluctuations occur in manufacturing processes in conjunction with the finer miniaturization of layout data, wherefore the precision of extracted LCR values can be enhanced by subjecting the LCR values found from layout data to corrections corresponding to those pattern width fluctuations. Accordingly, signal path propagation delay times can be found more accurately using those corrected LCR values, and suitable logic simulations can be performed."}
{"text": "The present application relates to devices for removing, capturing and rendering harmless bio-contaminants from air streams.\nFilters have long been used to remove contaminants in air handling systems. There have also been attempts to use ultra violet or UV lamps to destroy bacteria on a variety of surfaces, including filter surfaces. There is a long felt but unmet need for an effective purification unit that combines the bacteria destroying affects of UV radiation with effective purification by filtering in a more efficient and effective manner and simultaneously killing bacteria in the air stream and on the filtering medium.\nThe present invention provides a unique, simple combination filter and ultra-violet lamp assembly that includes a high purification filter and high power UV radiation lamp in a single replaceable and disposable unit capable of being easily installed into an existing air handling system, designed into a new air handing unit or for use in a stand alone room air purifier. The invention includes a tubular filter medium, preferably a pleated HEPA or ULPA filter mounted by opposing supporting end plates with one plate adapted for mounting a UV lamp and fixture in air sealed relationship to the filter medium such that the lamp is inside and surrounded by the filter medium and the power connecting portion of the fixture is outside the supporting end plate. The other supporting end plate has an air ingress opening therethrough.\nAccording to one embodiment, a self contained air purification unit is provided for replaceably mounting the filter and lamp assembly within an air handling system. The unit includes a housing having at least one removable wall for access to the filter and lamp assembly for installation and removal of the assembly. The housing is constructed to channel air flow through the housing for passage into the filter and lamp assembly and out through the filter medium. According to one aspect of the invention, a fan and motor assembly is mounted within the housing to draw air through the filter and in another the fan and motor assembly is mounted within the housing to push air through the filter."}
{"text": "Modern operating system kernels provide a mechanism called “asynchronous input/output (I/O)” where an execution thread submits an I/O request and then continues its execution while the I/O is performed in the background. Eventually the submitting thread (or perhaps another thread) checks or waits for the I/O to complete, and continues processing with the result of the operation. This mechanism avoids the inherent delays accompanying I/O requests that are serviced by long-latency devices, such as rotating disks and block storage devices. Typically such devices are mechanical in nature and must physically move to seek a track to read or write, which is orders of magnitude slower than the switching of electric current done by a processor in instruction-processing cycles.\nIf the storage system is simple (e.g., a block storage device) asynchronous I/O is simple to implement. However, if the storage system is complex (e.g., a file system), where a single I/O request from a thread translates to multiple interdependent requests to the device, implementation of asynchronous I/O becomes more difficult.\nCurrently, there are multiple ways to implement asynchronous I/O. One technique involves a state machine, where the kernel constructs a structure to track the state of the request, and submits the first of a series of I/O requests to the device. Each time a request completes, the kernel updates the structure and submits the next request. Eventually processing completes and the kernel signals the completion to the requesting thread. The drawback of this approach is that it is complicated to implement and hard to maintain.\nAnother technique involves passing the request to another thread, known as a helper thread. The requesting thread can then continue processing, while the helper thread processes the I/O request synchronously in the background. When the helper thread completes, it signals the requesting thread so that the requesting thread can then complete the I/O. While simple to implement, this approach incurs additional context switches, even in the case where the I/O could, in fact, be submitted immediately due to the needed data already being cached. These drawbacks reduce, and perhaps even nullify, the performance advantage gained by implementing asynchronous I/O in the first place.\nYet another technique involves running the process, submitting the first interdependent I/O in the series, and returning without waiting for the request to complete. When the I/O completes, the request is restarted from the beginning, so that the contents of the first I/O request are in the cache. Then, the thread proceeds to the second interdependent I/O in the series, submits the request, returns, and restarts from beginning when I/O completes, and so on. This approach avoids context switches, but is complex and requires extra work when restarting the process."}
{"text": "1. Field of the Invention\nThe present invention relates to an abnormality determination apparatus for an internal combustion engine, and particularly to a technique of determining that a passage for allowing exhaust gas discharged from a cylinder to be returned to the cylinder has been clogged.\n2. Description of the Background Art\nAn internal combustion engine provided with an EGR (Engine Gas Recirculation) system is known. With the EGR system, pumping loss and unburned gas can be reduced by means of an external EGR for allowing exhaust gas discharged from a cylinder to an exhaust passage to be recirculated to an intake passage.\nThe amount of exhaust gas recirculated by the EGR system from the exhaust passage to the intake passage is regulated by an EGR valve provided on an EGR gas passage. When the EGR valve opens, recirculation from the exhaust passage to the intake passage is performed. When the EGR valve closes, recirculation from the exhaust passage to the intake passage is stopped.\nWhen exhaust gas is recirculated into the cylinder, the combustion temperature within the cylinder decreases. Consequently, knocking hardly occurs. Since knocking hardly occurs, the need to retard ignition timing decreases. Thus, with recirculation of exhaust gas into the cylinder, ignition timing is advanced.\nAt this time, for example, if a passage provided for any of a plurality of cylinders for allowing exhaust gas to be returned into the cylinder is clogged, the advancement of ignition timing could be accompanied by occurrence of knocking. If the ignition timing is retarded in order to reduce knocking, a combustion state could be degraded in cylinders to which exhaust gas is recirculated normally. Therefore, it is necessary to identify the cylinder for which the passage for allowing exhaust gas to be returned into the cylinder is clogged, and to take appropriate measures.\nJapanese Patent Laying-Open No. 2010-25059 discloses, in paragraph 34 for example, that when an intake pressure sensor detects a decrease of the intake pressure below a target pressure, and a knock sensor detects occurrence of knocking at any particular cylinder, ignition timing of the particular cylinder is retarded by a spark plug without changing an opening position of an EGR valve.\nHowever, the intake pressure may decrease and knocking may occur even if a passage for allowing exhaust gas to be returned into the cylinder is not clogged. Therefore, with the method described in Japanese Patent Laying-Open No. 2010-25059, there could be a case where clogging of a passage for allowing exhaust gas to be returned into a cylinder is erroneously detected."}
{"text": "1. Field of the Invention\nThe present invention relates to a drive control technique for an optical modulator used in optical communication. In particular, the invention relates to a drive control apparatus and drive control method, which can stably control the drive of an optical modulator used in ultra high speed optical transmission and the like.\n2. Description of the Related Art\nRecently, in optical transmission systems, with the progress of wavelength division multiplexing (WDM) technology and higher transmission speed per wavelength, systems enabling of large capacity optical transmission are being realized. Also, for the modulation code form for transmission used in such systems, in order to enlarge the transmission capacity and transmission distance, forms such as optical duo binary code and carrier suppressed RZ (CS-RZ) code have been proposed, instead of conventional NRZ (non return to zero) code or RZ (return to zero) code (refer to Japanese Unexamined Patent Publication No. 2001-119344). In particular, the CS-RZ code is said to be the most promising modulation code form for the next generation, since it enables enlargement of transmission distance and transmission capacity by a simple method.\nAlso, in optical modulation for ultra high speed transmission, in order to suppress wavelength fluctuations (chirping) during modulation, an external modulation method has been used, for allowing a semiconductor laser to continuously emit light, to switch the light on and off by an external modulator. For the external modulator, a Mach-Zehnder optical modulator (to be referred to as MZ modulator hereunder) is prevailing.\nFIG. 31 shows an example of optical output characteristics of an MZ optical modulator with respect to a drive signal. In the conventional NRZ modulation method, on/off modulation of the optical output is performed by matching a high level and a low level of a drive signal, whose level fluctuates at a period T0, with an apex A of optical emission and an apex B of optical extinction in the optical output characteristics. In the following explanation, a difference of bias voltages corresponding to the apexes A and B of emission/extinction in the optical output characteristics which fluctuate periodically is represented by Vπ. According to this, an amplitude of the drive signal in the NRZ modulation method described above becomes Vπ.\nFor the MZ optical modulator as described above, the advantage is less chirping. However, there is a problem in that the optical output characteristics with respect to the drive signal drift timewise due to temperature change and a change with time, resulting in that interference between the codes occurs in the on/off level of the optical output. In order to solve this problem to stably control an operating point of the MZ optical modulator, a control at the operating point is necessary, such that, as shown in FIG. 31, in the case where a curve “a” showing the optical output characteristics with respect to the drive voltage moves to the left to become a curve “b”, a bias voltage by the drive signal is made lower according to the such a change, and in the case where the curve “a” moves to the right to become a curve “c”, the bias voltage is made higher according to such a change. In the conventional modulation method using the NRZ code or the RZ code, a compensation technique has been proposed in Japanese Patent Publication No. 3-251815, for superimposing a low frequency signal onto the drive signal to detect an fluctuation amount and a fluctuation direction of the operating point, and then controlling the bias voltage by feedback to hold the operating point normal.\nUnder the background as described above, also for the new modulation methods such as the optical duo binary modulation method and the CS-RZ modulation method and the like, an operating point control technique for the MZ optical modulator which can correspond to these modulation methods is required.\nAs an operating point control technique for an MZ optical modulator corresponding to the CS-RZ modulation method, there is the one disclosed in Japanese Unexamined Patent Publication No. 2001-119344. In this conventional technique, as shown in FIG. 32, a transmission section of the CS-RZ modulation method comprises: a transmission light source 110 that generates continuous light; an MZ optical modulator 120 on a data modulation side, that is driven in accordance with a drive signal obtained by adjusting, by a drive circuit 121, a data signal corresponding to the NRZ modulation method so that the amplitude thereof becomes Vπ, and an MZ optical modulator 130 on a clock modulation side, that is driven in accordance with a drive signal obtained by adjusting, by a drive circuit 131, a clock signal having a frequency of ½ the bit rate of the data signal so that the amplitude thereof becomes 2Vπ.\nIn such a transmission section of the CS-RZ modulation method, by adjusting a crossing point in an output from the optical modulator 120 on the data modulation side, to an optical extinction phase of the optical modulator 130 on the clock modulation side, an optical signal having the same waveform as for when the RZ code is used is output as an optical output. Therefore, for the optical modulator 120 on the data modulation side, since the drive signal of the NRZ code form with amplitude Vπ is used in the same manner as for the aforementioned case shown in FIG. 31, the operating point can be optimized by applying the conventional operating point control method. On the other hand, for the optical modulator 130 on the clock modulation side, for example as shown in FIG. 33, since the waveform of the drive signal is a sine wave with amplitude 2Vπ, it becomes necessary to control the operating point so that the amplitude of the drive signal at the time of optical extinction becomes zero.\nIn addition, for the operating point control method of the MZ optical modulator corresponding to the optical duo binary modulation method, the technique disclosed in Japanese Unexamined Patent Publication No. 2000-162563 is known. In this conventional technique, for example as shown in FIG. 34, different to the case of the modulation method using the NRZ code or the RZ code, a drive signal corresponding to the optical duo binary code, being a differential code with amplitude 2Vπ, is supplied to an MZ optical modulator 210. This drive signal is superimposed with an operating point control signal of low frequency generated by an oscillator 212. A part of the optical output from the MZ optical modulator 210 is branched by an optical coupler 220, and then sent to a phase comparator 223 via an optical receiver 221 and an I/V converter 222. In the phase comparator 223, a phase comparison is performed on the superimposed component contained in the optical output, and the operating point control signal, to thereby detect the fluctuation amount and the fluctuation direction of the operating point. Then, the operating point is optimized by feeding back the result of the phase comparison to a control of a bias supply circuit 225. Furthermore, in the conventional technique described above, for example as shown in FIG. 35, a compensation method has also been proposed, for performing a control such that the frequency of the superimposed component contained in the optical output becomes a maximum, by utilizing the phenomenon where the frequency of the superimposed component becomes twice the frequency of the operating point control signal when the bias voltage coincides with the operating point, to hold the operating point normal.\nConsidering the control of the operating point of the MZ optical modulator on the clock modulation side in the CS-RZ modulation method by applying the conventional operating point control technique as described above, even if an operating point control signal with a low frequency is superimposed on the clock signal, since the waveform of the clock signal is a sine wave, the optical output from the optical modulator only becomes a “1” level (apex A of optical emission) temporally, and there is no point where a “1” level is continuous for the optical output, as in the case where an operating point control signal is superimposed on a rectangular wave data signal. Therefore, the superimposed component with a low frequency corresponding to the operating point control signal rarely appears in the optical output from the optical modulator. Consequently, in the conventional method of superimposing the operating point control signal, it is difficult to detect the fluctuation amount and the fluctuation direction of the operating point in the MZ optical modulator on the clock modulation side.\nAlso, in the method which controls the operating point by detecting that the frequency of the superimposed component contained in the output from the optical modulator becomes twice the frequency of the operating point control signal, in the MZ light modulator-on the clock modulation side, since the point where the optical output becomes a maximum is only temporary as described above, the detection signal becomes faint so that it is difficult to determine the frequency of the superimposed component. Moreover, in the conventional method, there is also a problem in that when the amplitude of the drive signal is shifted from 2Vπ, then in the optical output, two bias points exist which generate a superimposed component with twice the frequency of the operating point control signal, and in the case where the operating point is controlled based on either of the bias points as a reference, only output waveforms with deviated duty are obtained, so that normal modulation operation cannot be realized.\nFurthermore, in the transmission section of the CS-RZ modulation method as shown in FIG. 32, in order to obtain a desired optical output, an adjustment or a control is required at several points such as the amplitude and the phase between the two drive signals of differential form corresponding to the clock signal, and the phases between the clock signal and the data signal. Therefore, there is a disadvantage in that when providing the additional configurations for the conventional drive control, this invites a higher cost and a larger size for the transmission section."}
{"text": "Cancer is the second leading cause of death in the United States, exceeded only by heart disease. (Cancer Facts and Figures 2004, American Cancer Society, Inc.). Despite recent advances in cancer diagnosis and treatment, surgery and radiotherapy may be curative if a cancer is found early, but current drug therapies for metastatic disease are mostly palliative and seldom offer a long-term cure. Even with new chemotherapies entering the market, the need continues for new drugs effective in monotherapy or in combination with existing agents as first line therapy, and as second and third line therapies in treatment of resistant tumors.\nCancer cells are by definition heterogeneous. For example, within a single tissue or cell type, multiple mutational “mechanisms” may lead to the development of cancer. As such, heterogeneity frequently exists between cancer cells taken from tumors of the same tissue and same type that have originated in different individuals. Frequently observed mutational “mechanisms” associated with some cancers may differ between one tissue type and another (e.g., frequently observed mutational “mechanisms” leading to colon cancer may differ from frequently observed “mechanisms” leading to leukemia). It is therefore often difficult to predict whether a particular cancer will respond to a particular chemotherapeutic agent (Cancer Medicine, 5th edition, Bast et al., B. C. Decker Inc., Hamilton, Ontario).\nComponents of cellular signal transduction pathways that regulate the growth and differentiation of normal cells can, when dysregulated, lead to the development of cellular proliferative disorders and cancer. Mutations in cellular signaling proteins may cause such proteins to become expressed or activated at inappropriate levels or at inappropriate times during the cell cycle, which in turn may lead to uncontrolled cellular growth or changes in cell-cell attachment properties. For example, dysregulation of receptor tyrosine kinases by mutation, gene rearrangement, gene amplification, and overexpression of both receptor and ligand has been implicated in the development and progression of human cancers."}
{"text": "1. Field of the Invention\nThe present invention relates to a surface emitting laser, a surface emitting laser array, and an optical apparatus having the surface emitting laser array.\n2. Description of the Related Art\nA vertical-cavity surface-emitting laser (VCSEL), by which it is easy to form a two-dimensional array as an exposure light source for forming a latent image on a photoreceptor as a scanning surface, is used for a scanning optical apparatus in an electrophotographic image forming apparatus for use in a laser beam printer or a digital copier.\nAn example of an optical system of an image forming apparatus using a surface emitting laser is disclosed in Japanese Patent Application Laid-Open No. 2008-052197. According to Japanese Patent Application Laid-Open No. 2008-052197, a surface emitting laser having a plurality of light emitting points is used as a multi-beam laser light source. In order to limit the light beam width of a light beam emitted from the light source and to form the beam into a desired beam shape, the optical system has an aperture stop common to each laser light source. The formed beams are deflected and scanned by an optical deflector and imaged in a spot shape on a photoreceptor drum surface as a scanning surface by an imaging optical system."}
{"text": "Fluid supply networks, such as municipal water supply systems, typically test the flow and pressure characteristics of outlets, such as hydrants, on the system to determine the capabilities of the system to deliver an adequate flow of water during high water use events such as the fighting of a fire using water from the supply system.\nNot only are these characteristics of the municipal water supply system a concern of the municipality in gauging the likely ability of the fire department to fight fires, but insurers of property are also concerned with this capacity in setting insurance rates for the property. Accurate and current data for each hydrant indicates how much water is available from each hydrant, which may affect fire suppression practices at different points on the system so as not to place greater fire fighting water demands on the hydrant than the individual hydrant or hydrants can supply. The hydrant flow characteristics are not measured in isolation from other proximate hydrants, so that fire fighting personnel are aware whether the system will support the use of multiple proximate hydrants without simply diverting the flow from hydrants already opened and being utilized. Additionally, the use of fire fighting equipment to draw water from the hydrant using a pump can increase the flow capacity of water from the water supply system, but attempting to pump water from the hydrant in excess of the safe capacity of water available from the hydrant at that point in the supply system can have effects detrimental to the system due to the creation of a negative pressure condition created in the main. These detrimental effects include triggering a collapse of the water main or the introduction of contaminants into the main through cracks or joints in the pipe. Thus, knowledge of the actual flow capabilities of the water supply system at a hydrant is needed as guidance to fire fighting personnel seeking to draw water from the hydrant in emergencies.\nSimilarly, in an actual firefighting situation, it is desirable to monitor the residual pressure in the water supply system as water is being pumped or drawn out of the supply system by fire fighting equipment to ensure that the pressure in the supply system does not fall below a desired threshold pressure and possibly cause the aforementioned problems.\nFurther, flow test data can provide information about the water supply system so that system managers can estimate the capabilities of water mains, such as flow rates, and plan system upgrades and expansions or determine the resistance to water or fluid flow in a fluid main. The latter resistance and flow velocity may be used by engineers, such as by coefficient factor (C-factor) testing, to determine if a main is plugging or becoming clogged. In this latter testing, flow velocity (V) is related to a flow coefficient (C) by the Hazen-Williams flow formula:\n      V    =          1.318      ⁢                        C          ⁡                      (                          R              H                        )                          0.63            ⁢                        (          S          )                0.54                  Where:                                V          =                    ⁢                      flow velocity, ft/s                                                        C          =                    ⁢                      flow coefficient                                                                    R            H                    =                    ⁢                      hydraulic radius, ft                                                                  ⁢                                                    Note:                             ⁢                              R                H                                      =                                          1                /                4                            ⁢                              (                Di                )                            ⁢                                                          ⁢                              for pipe flowing full or half full                                                                                   S          =                    ⁢                      hydraulic slope, ft/ft                                                                    D            i                    =                    ⁢                      pipe inside diameter, ft                              Also,\n      Q    =          0.006756      ⁢                          ⁢              CD        i        2.63            ⁢              f        0.54                  Where:                                Q          =                    ⁢                      flow rate, gal/min                                                        C          =                    ⁢                      flow coefficient                                                                    D            i                    =                    ⁢                      pipe inside diameter, in                                                                    f            =                        ⁢                                          head loss, ft of                            ⁢                                                          ⁢                              H                2                            ⁢              O              ⁢                              /1000 ft                                              ⁢                                          ⁢                      f            =                                                            0.2083                  ⁡                                      [                                          100                      C                                        ]                                                  1.85                            ⁢                                                Q                  1.85                                                  D                  i                  4.86                                                              ⁢                                          ⁢                      Where:                                                        f          =                    ⁢                                    friction loss, ft of                         ⁢                                                  ⁢                          H              2                        ⁢            O            ⁢                          /1000 ft                                                                    C          =                    ⁢                      flow coefficient                                                        Q          =                    ⁢                      flow rate, gal/min                                                                    D            i                    =                    ⁢                      pipe inside diameter, in                              \nFlow rate (e.g., gal/minutes) from a hydrant is related to the flow coefficient C and pressures in the systems as follows:\n      Q    =          0.442      ⁢              D        i        2.63            ⁢                        C          ⁡                      [                                                            P                  1                                -                                  P                  2                                            L                        ]                          0.54                  Where:                                Q          =                    ⁢                      flow rate, gal/min                                                                    D            i                    =                    ⁢                      pipe inside diameter, in                                                        C          =                    ⁢                      flow coefficient                                                                    P            1                    ,                                    P              2                        =                        ⁢                          gauge pressures, psi                                                                    L          =                    ⁢                      pipe length, ft                              \nHydrant flow characteristics can affect decisions as to what fire protection and fire resistance features are required for areas of new developments, and where priorities should be placed with respect to upgrading older, smaller water mains. Such testing can indicate systemic weaknesses such as main clogging, failing water mains such as when residual pressures imposed upon water mains become too high by pulling water from the main at a flow hydrant, and compromised valves.\nTesting is typically conducted on a periodic basis, since water supply systems are constantly being affected by changing conditions, including improvements to the system, deterioration of parts of the system, and changes in usage of the system, etc. The testing of hydrants in a water supply system may follow the requirements of National Fire Protection Association (NFPA) No. 291, entitled “First Flow Testing and Marking of Hydrants”. In general, the testing of hydrants in a municipal water supply system where the testing is directed to flow rates and the effect thereof. These tests involve the measurement of static pressure, residual pressure and pitot pressure with respect to the subject hydrant being tested (a test hydrant). C-factor testing may only involve measuring or determining flow velocity simultaneously with residual pressure and calculating resistance to fluid flow using other known characteristics of the system. However, the procedures of both of these tests are not performed only at the subject hydrant being tested, but also at another hydrant on the same water supply system as the subject test hydrant. This testing provides an accurate idea of the subject hydrant's individual characteristics and efficiency as well as other characteristics of the system such as a flow coefficient. More specifically, as to flow characteristics which involve the effect of flow rate on a system, static pressure and then residual pressure are measured at the subject test hydrant, while pitot pressure is measured in a flow of water at a proximate location on the water supply system, such as a fire hydrant (the “flow hydrant”) downstream or adjacent to the subject hydrant being tested on the water supply system (adjacent especially in the case of determining the effect of flow rate on residual pressure, or other access points to the water flow, even from a nearby residential water service). In the latter circumstance when determining the effect of flow rate on residual pressure, because the water in the supply system is flowing at the proximate or adjacent flow hydrant, measurements taken at the subject test hydrant are substantially isolated from effects such as friction loss, and the measurements of these characteristics is thus more accurate.\nIt will be noted that although the hydrants are proximate or adjacent to each other in the water supply system, the hydrants are typically widely separated in a geographic sense, and personnel are usually stationed at each of the hydrants to conduct the testing.\nJust as significant as the physical separation of the subject test hydrant and the proximate flow hydrant is the temporal requirements of the testing. More specifically, while the static pressure is typically measured at the subject test hydrant in testing involving flow rate and the effects thereof just prior to opening the proximate flow hydrant to flow water from the system, the measurement of pitot pressure at the proximate flow hydrant and residual pressure at the subject test hydrant must occur simultaneously or substantially simultaneously. Substantially simultaneously means that errors caused by not precisely temporally linking the determined residual pressure with a determined flow rate using pitot pressure should not exceed 10% and preferably not more than 5% and even more preferably not more than 1%.\nIn the past systems, methods and apparatuses used in flow rate, flow velocity or C-factor testing, one person alone has not been able to accurately conduct these tests because of the geographical and temporal requirements of these tests. Whether involving flow rate testing or flow velocity used in C-factor testing, in the past at least two persons have been required to take action at each hydrant, the test hydrant and flow hydrant. Thus, the expense of the testing process is increased by the personnel costs. Also, the simultaneous timing of the measurements has not been always reliable, as clocks must be synchronized and/or the testing personnel must communicate the exact time of the taking of the measurement, once the proximate hydrant valve has been fully opened and the pressure in the water main has stabilized. Communication of this timing is thus often performed verbally by the personnel over a portable radio system. Also, there is rarely, if ever, any independent verification of the simultaneous timing of the taking of these measurements.\nTesting equipment which monitors pressure at fire hydrants are known where a hydrant pressure measuring device and recorder are attached to a hydrant being tested for static and residual pressure as well as attached to a hydrant which is permitted to flow and flow rate is determined. In these systems a hand held computer or personal digital assistant is configured to download pressure or flow data being measured at each hydrant. Clocks in each unit on each hydrant are carefully synchronized so that data being measured at each hydrant can be correlated and the effect of fluid flow at a flow rate at a flow hydrant can be seen on a residual pressure on a test hydrant. There is no communication between the devices on each hydrant to assure that the collected data concerning a pitot pressure and/or flow rate will directly correlate to a residual pressure being measured at a test hydrant. Moreover, for these types of devices to work, their clocks at each hydrant have to be carefully and precisely synchronized.\nFurther increasing the expense and complication of the testing are situations where more than one hydrant needs to be opened to achieve a desirable drop in residual pressure approaching 25 percent from the static pressure. This adds to the personnel expense and complication to the timing of the opening of the hydrants.\nIn the fire-fighting situation, where it is desirable to monitor the pressure in the water supply system and the effect on such pressure of a flow rate from a flow hydrant, an additional fire fighter or municipal employee must be stationed at a hydrant proximate to the flow hydrant from which water is being drawn in order to monitor the residual pressure in the supply system. This situation thus also requires additional personnel simply to monitor flow conditions in the supply system. In fire fighting situations, the term “fire flow” is sometimes used. As discussed herein, fire flow means the rate at which water may be drawn from the system (such as gallons per minute) without the residual pressure dropping below 20 psi.\nIt is also desirable to coordinate the opening of the proximate hydrant (or the downstream hydrant in the case of C-factor testing) and the taking of pressure and flow readings in order to minimize the time that the flow hydrant is opened. This minimizes the amount of water that flows from the proximate or downstream flow hydrant, and thereby minimizes the amount of water that is wasted and that needs to be disposed of.\nUsing the static, pitot, and residual pressure measurements, the flow rate in gallons per minute may be calculated using the formula:\n  Q  =      29.83    ⁢                  ⁢          cd      2        ⁢          p      \nwhere Q=observed flow, c=coefficient, d=outlet diameter, p=pitot pressure.\nThe available flow may be calculated, subject to some qualifications, using the formula:\n      Q    R    =            Q      F        ×                  h        r        0.54                    h        f        0.54                            where Qr=observed flow, hr is the drop in pressure from the static pressure to the desired residual baseline and hf is the drop in psi from static pressure to the actual residual pressure that was measured during the test.        \nDue to the accuracy required in making the measurements, and the labor intensive nature of the taking of the measurements over sometimes long distances, it is believed that there is needed a system for administering and recording data relative to flow testing for the effect of flow rate on residual pressure or for C-factor testing involving flow velocity and conducting these tests in a manner that is able to increase the accuracy of and decrease the personnel needed for such tests. These measurements not only include measurements of static and residual pressures, determination of flow rates and flow velocities such as from measuring pitot pressure, but also include identifying specific geographical locations of test and flow hydrants using a global positioning system (GPS), correlating the data to record the location of each hydrant into a memory or storage device and associating the location of each hydrant with the data taken at each hydrant. Hence, there is not only a need for making measurements concerning hydrants and valves with devices which “talk” to each other via wireless communications, but also geographically map the hydrants or valve system and associate the measurements taken as to the hydrant or valve with a position on the geographical map."}
{"text": "The present invention relates to a polar transmitter, and more particularly, to an amplitude modulation circuit in a polar transmitter and a method for calibrating the amplitude offset in the polar transmitter.\nConventional polar transmitters have already been disclosed and discussed in various literatures such as U.S. Patent Publication No. 20060089111, whose content is incorporated herein by reference. Furthermore, in transmission systems in which a combined AM and PM signal is used in a power control loop having a wide range of power output levels and to meet strict GSM/enhanced data rates for GSM evolution (EDGE) spectral emissions limitations, a high level of precision is typically required when converting the amplitude modulated portion of the transmit signal from the digital domain to the analog domain. This typically requires a DAC having 11-bit resolution, and when a larger cover range for the DC offset is required in U.S. Patent Publication No. 20060089111, it has to use a DAC having higher resolution such as 12-bit or 13-bit. However, such a DAC is costly to implement, both with respect to area on the circuit and power consumption. For example, the 12-bit DAC has twice larger area than the 11-bit DAC."}
{"text": "1. Field of the Disclosure\nEmbodiments disclosed herein relate generally to the processing of a C3 to C6 hydrocarbon cut from a cracking process, such as steam or fluid catalytic cracking, primarily for conversion of C4 olefins to propylene via metathesis.\n2. Background\nIn typical olefin plants, such as illustrated in U.S. Pat. No. 7,223,895, there is a front-end demethanizer for the removal of methane and hydrogen followed by a deethanizer for the removal of ethane, ethylene and C2 acetylene. The bottoms from this deethanizer tower consist of a mixture of compounds ranging in carbon number from C3 to C6. This mixture may be separated into different carbon numbers, typically by fractionation.\nThe C3 cut, primarily propylene, is removed as product and is ultimately used for the production of polypropylene or for chemical synthesis such as propylene oxide, cumene, or acrylonitrile. The methyl acetylene and propadiene (MAPD) impurities must be removed either by fractionation or hydrogenation. Hydrogenation is preferred since some of these highly unsaturated C3 compounds end up as propylene thereby increasing the yield.\nThe C4 cut, consisting of C4 acetylenes, butadiene, iso- and normal butenes, and iso- and normal butane can be processed in many ways. A typical steam cracker C4 cut contains the following components in weight %:\nTABLE 1Typical C4 cut components and weight percentages.C4 AcetylenesTraceButadiene33%1-butene15%2-butene 9%Isobutylene30%Iso- and Normal Butanes13%The components in a refinery or FCC based C4 cut are similar, with the exception that the percentage of paraffins are considerably greater.\nTypically, the butadiene and C4 acetylenes are removed first. This can be accomplished by either hydrogenation or extraction. The product from butadiene and C4 acetylene removal is designated Raffinate I. If extraction is employed, the remaining 1-butene and 2-butene remain essentially in the same ratio as that of the initial feedstock. If hydrogenation is employed, the initial product from butadiene hydrogenation is 1-butene. Subsequently, hydroisomerization occurs within the same reaction system changing the 1-butene to 2-butene. The extent of this reaction depends upon catalyst and reaction conditions within the hydrogenation system. However, it is common practice to limit the extent of hydroisomerization in order to avoid “over hydrogenation” and the production of butanes from butenes. This would represent a loss of butene feedstock for downstream operations. The butenes remaining in the mixture consist of normal olefins (1-butene, 2-butene) and iso-olefins (isobutylene). The balance of the mixture consists of both iso- and normal-butanes from the original feed plus what was produced in the hydrogenation steps and any small quantity of unconverted or unrecovered butadiene.\nA Raffinate I stream can be further processed in many ways. A Raffinate II stream is by definition a stream following isobutylene removal. Isobutylene can be removed in a number of ways. It can be removed via fractionation. In fractionation isobutane will be removed along with the isobutylene. In addition, some fraction of the 1-butene will be lost as well. The resultant Raffinate II will contain primarily normal olefins and paraffins and minimal isoolefins and isoparaffins. Isobutylene can also be removed via reaction. Reactions include: reaction with methanol to form MTBE, reaction with water to form tertiary butyl alcohol, or reaction with itself to form a C8 gasoline component. In all reaction cases, the paraffins are not removed, and thus the mixture will contain both normal and iso-paraffins. The paraffin content and composition of the Raffinate II impacts downstream processing options.\nThe butenes have many uses. One such use is for the production of propylene via metathesis. Another is for the production of ethylene and hexene via metathesis. Conventional metathesis involves the reaction of normal butenes (both 1-butene and 2-butene) with ethylene (principally the reaction of 2-butene with ethylene to form propylene). These reactions occur in the presence of a group VIA or VIIA metal oxide catalyst, either supported or unsupported. The paraffin components of the reaction feed are essentially inert and do not react, and are typically removed from the process via a purge stream in the separation system that follows the metathesis reactor. Typical catalysts for metathesis are tungsten oxide supported on silica or rhenium oxide supported on alumina. Examples of catalysts suitable for the metathesis of olefins are described in U.S. Pat. Nos. 6,683,019 and 6,420,619, for example. Isobutylene (isobutene) may be removed from the feedstock prior to the metathesis reaction step. The reaction of isobutylene with ethylene is non-productive and reaction with itself and/or other C4's is limited in the presence of excess ethylene. Non-productive reactions essentially occupy catalyst sites but produce no product. If allowed to remain in the feed to the metathesis unit, the concentration of this non-reactive species would build up creating capacity limitations. The reaction of 1-butene with ethylene is also non-productive. However, it is common to employ a double bond isomerization catalyst within the metathesis reactor to shift 1-butene to 2-butene and allow for continued reaction. Typical double bond isomerization catalysts include basic metal oxides (Group IIA), either supported or unsupported. Magnesium oxide and calcium oxide are examples of such double bond isomerization catalysts that may be physically admixed with the metathesis catalyst. No equivalent co-catalyst exists for the skeletal isomerization of isobutylene to normal butene. In the case of a conventional metathesis system employing both a metathesis catalyst and a co-mixed double bond isomerization catalyst, the butadiene must be removed to a level of less than 500 ppm to avoid rapid fouling of the double bond isomerization catalyst. The metathesis catalyst itself can tolerate butadiene levels up to 10,000 ppm.\nIn some cases, an isobutylene removal step is employed prior to metathesis. Options include reacting it with methanol to produce methyl tertiary butyl ether (MTBE) or separating the isobutylene from the butenes by fractionation. U.S. Pat. No. 6,358,482 discloses the removal of isobutylene from the C4 mixture prior to metathesis. This scheme is further reflected in U.S. Pat. Nos. 6,075,173 and 5,898,091. U.S. Pat. No. 6,580,009 discloses a process for the production of propylene and hexene from a limited ethylene fraction. For molar ratios of ethylene to butenes (expressed as n-butenes) from 0.05 to 0.60, the inventors utilize a Raffinate II stream as the C4 feedstock.\nThe typical metathesis process takes the Raffinate I feedstock and removes the majority of the isobutylene via fractionation, as described above to form a Raffinate II. In this step, the isobutene is removed as well plus some quantities of normal butenes, dependent upon the fractionation conditions. The Raffinate II is then admixed with ethylene, passed through guard beds to remove poisons, vaporized and preheated and fed to the metathesis reactors. The operating conditions are typically 300° C. and 20 to 30 bar pressure. The reactor effluent following heat recovery is then separated in a fractionation system. First the ethylene is recovered overhead in a first tower and recycled to the reactor system. The tower bottoms are then sent to a second tower where the propylene is recovered overhead. A side draw is taken containing the majority of the unconverted C4 components and recycled to the reactor. The tower bottoms containing the C5 and heavier products plus C4 olefins and paraffins are sent to purge. The purge rate is typically fixed to contain sufficient C4 paraffins to avoid their buildup in the reactor recycle stream. In some cases, a third tower is employed on the tower bottoms stream to separate the C4 components overhead and the C5 and heavier components as a bottoms stream.\nU.S. Pat. No. 6,271,430 discloses a two-step process for the production of propylene. The first step consists of reacting 1-butene and 2-butene in a raffinate II stream in an auto-metathesis reaction to form propylene and 2-pentene. The products are then separated in the second step. The third step reacts specifically the 2-pentene with ethylene to form propylene and 1-butene. This process utilizes the isobutylene free raffinate II stream. The pentenes recycled and reacted with ethylene are normal pentenes (2-pentene).\nIsobutylene removal from the C4 stream can also be accomplished by employing a combined catalytic distillation hydroisomerization deisobutyleneizer system to both remove the isobutylene and recover n-butenes at high efficiency by isomerizing the 1-butene to 2-butene with known isomerization catalysts and thus increasing the volatility difference. This technology combines conventional fractionation for isobutylene removal with hydroisomerization within a catalytic distillation tower. In U.S. Pat. No. 5,087,780 to Arganbright, 2-butene is hydroisomerized to 1-butene as the fractionation occurs. This allows greater than equilibrium amounts of 1-butene to be formed as the mixture is separated. Similarly, 1-butene can be hydroisomerized to 2-butene in a catalytic distillation tower. In separating a C4 stream containing isobutylene, 1-butene, and 2-butene (plus paraffins), it is difficult to separate isobutylene from 1-butene since their boiling points are very close. By employing simultaneous hydroisomerization of the 1-butene to 2-butene with fractionation of isobutylene, isobutylene can be separated from the normal butenes at high efficiency.\nThe metathesis reaction described above is equimolar, i.e., one mole of ethylene reacts with 1 mole of 2-butene to produce 2 moles of propylene. However, commercially, in many cases, the quantity of ethylene available is limited with respect to the quantity of butenes available. In addition, the ethylene is an expensive feedstock and it is desired to limit the quantities of ethylene used. As the ratio of ethylene to butenes is decreased, there is a greater tendency for the butenes to react with themselves which reduces the overall selectivity to propylene.\nThe metathesis catalysts and the double bond isomerization catalysts are quite sensitive to poisons. Poisons include water, CO2, oxygenates (such as MTBE), sulfur compounds, nitrogen compounds, and heavy metals. It is common practice to employ guard beds upstream of the metathesis reaction system to insure the removal of these poisons. It does not matter if these guard beds are directly before the metathesis reaction system or further upstream as long as the poisons are removed and no new poisons are subsequently introduced.\nMetathesis reactions are very sensitive to the location of the olefin double bond and the stereo-structure of the individual molecules. During the reaction, the double bond on each pair of olefins adsorb on the surface and exchange double bond positions with the carbon groups on either sides of the double bonds. Metathesis reactions can be classified as productive, half productive or non-productive. As described above, non-productive reactions result in essentially no reaction taking place. When the double bonds shift with metathesis reaction, the new molecules are the same as the originally adsorbed molecules thus no productive reaction occurs. This is typical for reactions between symmetric olefins or reactions between ethylene and alpha olefins. If fully productive reactions occur, new products are generated no matter which orientation the molecules occupy the sites. The reaction between ethylene and 2-butene to form two propylene molecules is a fully productive reaction. Half productive reactions are sterically inhibited. If the pair of olefins adsorb in one orientation (typically the cis position with respect to the attached R groups), when the double bonds shift, new products are formed. Alternately if they adsorb in a different steric configuration (the trans position), when the bonds shift, the identical olefins are formed and thus no new products are formed. The various metathesis reactions proceed at different rates (a fully productive reaction is usually faster than a half productive reaction). Table 2 summarizes the reactions between ethylene and various butenes and the reactions between the butenes themselves.\nThe reactions listed in Table 2 represent the base reaction with ethylene (reaction 1, 4 and 5) as well as the reactions between the various C4 olefins. It is especially important to make a distinction between the selectivity to propylene from total C4 olefins (including isobutylene) and the selectivity to propylene from the normal C4 olefins involved in the reaction. The reaction of isobutylene with 2-butene (reaction 6) produces propylene and a branched C5 molecule. For this reaction, propylene is produced at 50 molar % selectivity from total C4's (similar to reaction 2) but at a 100 molar % selectivity from the normal C4 (2-butene). For the purposes of definitions, conventional metathesis is defined as the reaction of the C4 olefin stream with ethylene. However, the C4 stream can also react in the absence of ethylene as a feedstock. This reaction is called auto or self metathesis. In this case, reactions 2, 3, 6, and 7 are the only possible reactions and will occur at rates dependent upon the feedstock composition.\nTABLE 2Molar %Molar %SelectivitySelectivity(C3H6 from(C3H6 fromNo.ReactionTypeRatetotal C4s)n-C4s)12-butene + ethylene →FullyFast1001002 propyleneProductive(Conventional Metathesis)21-butene + 2-butene →FullyFast5050Propylene + 2-penteneProductive31-butene + 1-butene →HalfSlow00Ethylene + 3-hexeneProductive4Isobutylene + Ethylene →Non-No——No reactionproductiveReaction51-butene + ethylene →Non-NoNo reactionproductiveReaction6Isobutylene + 2-butene →FullyFast50100Propylene + 2-methyl 2-buteneProductive7Isobutylene + 1-butene →HalfSlow00ethylene + 2-methyl 2 penteneproductive\nIn conventional metathesis, the focus is to maximize reaction 1 to produce propylene. This will maximize the selectivity to propylene. As such, excess ethylene is used to reduce the extent of the reactions of butenes with themselves (reactions 2, 3, 6, and 7). The theoretical ratio is 1/1 molar or 0.5 weight ratio of ethylene to n-butenes but it is common in conventional metathesis to employ significantly greater ratios, typically, 1.3 or larger molar ratio to minimize reactions 2, 3, 6 and 7. Under conditions of excess ethylene, and due to the fact that both isobutylene and 1-butene do not react with ethylene (see reactions 4 and 5), two process sequences are employed. First, the isobutylene is removed prior to metathesis. If isobutylene is not removed, it will build up as the n-butenes are recycled to achieve high yield. Second, 1-butene is isomerized to 2-butene by including a double bond isomerization catalyst such as magnesium oxide admixed with the metathesis catalyst. Note that this catalyst will not cause skeletal isomerization (isobutylene to normal butylenes) but only shift the double bond from the 1 position to the 2 position for the normal butenes. Thus by operating with excess ethylene, eliminating isobutylene from the metathesis feed prior to reaction, and employing a double bond isomerization catalyst, reaction 1 is maximized. Note, however, that by removing the isobutylene, potential production of propylene or other products is lost.\nWhen there is limited or no fresh ethylene (or excess butylenes for the ethylene available), there are currently two options available for propylene production. In these cases, the first option will first remove the isobutylene and then process the normal butenes with whatever ethylene is available. The entire n-butenes-only mixture is subjected to metathesis with the available ethylene. Ultimately, if there is no fresh ethylene available, the C4's react with themselves (auto metathesis). Under low ethylene conditions, reactions 2, 3, 6 and 7 will occur, all leading to a lower propylene selectivity (50% or lower versus 100% for reaction 1). The lower selectivity results in lower propylene production. Note that reactions 6 and 7 will be minimized as a result of the removal of isobutylene (to low levels but not necessarily zero). Alternately, the molar flows of ethylene and butenes can be matched by limiting the flow of butenes to produce conditions where there is a high selectivity of the normal butenes to propylene via reaction 1. By limiting the flow of n-butenes to match ethylene, the production of propylene is limited by the reduced butenes flow.\nPentenes and some hexenes are formed to some extent in the conventional metathesis case with low ethylene via reactions 2 and 3. The volume of these components will depend upon the ethylene/n-butenes ratio with a lower ratio producing more C5 and C6 components. In the conventional prior art case where isobutylene is removed before any metathesis, these C5 and C6 olefins are normal olefins since no skeletal isomerization occurs. It is possible to recycle these olefins back to the metathesis step where, for example, the reaction with ethylene and 2-pentene will occur yielding propylene and 1-butene. The 1-butene is recovered and recycled. Note however, with limited ethylene, reaction 1 can occur only to the limit of the ethylene availability. Ultimately these non-selective byproducts, pentenes and hexenes, must be purged from the system.\nU.S. Pat. No. 6,777,582 discloses a process for the auto-metathesis of olefins to produce propylene and hexene. Therein, auto-metathesis of a mixed normal butenes feed in the presence of a metathesis catalyst operates without any ethylene in the feed mix to the metathesis reactor. Some fraction of the 2-butene feed may be isomerized to 1-butene and the 1-butene formed plus the 1-butene in the feed react rapidly with the 2-butene to form propylene and 2-pentene. The feed to the reactor also includes the recycle of the 2-pentene formed in the reactor with unreacted butenes to simultaneously form additional propylene and hexene. The 3-hexene formed in the reaction may be isomerized to 1-hexene.\nIn U.S. Pat. No. 6,727,396, ethylene and hexene-1 are produced from butene-1 by metathesis of butene-1 and isomerization of the hexene-3 produced therein to hexene-1. The initial starting material is a mixed butene stream wherein butene-1 is isomerized to butene-2 with isobutylene being separated therefrom, followed by isomerization of butene-2 to butene-1, with the butene-1 being the feed to the metathesis.\nIn U.S. Pat. No. 7,214,841, the C4 cut from a hydrocarbon cracking process is first subjected to auto-metathesis prior to any isobutylene removal and without any ethylene addition, favoring the reactions which produce propylene and pentenes. The ethylene and propylene produced are then removed leaving a stream of the C.sub.4's and heavier components. The C.sub.5 and heavier components are then removed leaving a mixture of 1-butene, 2-butene, isobutylene, and iso- and normal butanes. The isobutylene is next removed preferably by a catalytic distillation hydroisomerization de-isobutyleneizer. The isobutylene-free C4 stream is then mixed with the product ethylene removed from the auto-metathesis product together with any fresh external ethylene needed and subjected to conventional metathesis producing additional propylene.\nAnother use of the C4 olefin stream is as a feedstock to an olefin cracking process, where the olefins are reacted with themselves over a zeolitic catalyst to produce a mixture comprising ethylene, propylene, and aromatics (such as benzene). Similar to the metathesis process, the paraffins are inert in this cracking process and must be removed from the process via a purge stream. U.S. Pat. No. 6,307,117 and U.S. Patent Application Publication No. 20050080307 both describe such a process. A mixture of typically C4 to C6 olefins as well as paraffins is vaporized and fed to a reactor filled with a crystalline zeolitic catalyst and operating at a temperature between 450 and 600° C. and a pressure between 10 and 70 psia. The reactor effluent is first sent to a compression step. The cracking reactor system operates at relatively low pressure to avoid fouling of the catalyst in the cracking reactor. In order to reduce the energy costs due to refrigeration in the subsequent separation system, the pressure is typically increased to pressures on the order of 12 to 25 barg. This allows the subsequent fractionating towers to utilize cooling water instead of refrigeration in the overhead condensation step. The compression effluent is then sent to a separation system where the ethylene and propylene are recovered along with an aromatics stream. The ethylene and propylene is recovered overhead in a first tower. Unlike metathesis, these products contain sufficient quantities of ethane and propane that additional purification of this stream is necessary. This can be accomplished by additional fractionation or by utilizing the recovery system of an adjacent facility such as an olefins plant. The tower bottoms contains C4, C5 and C6+ paraffins and aromatics. This is sent to a second tower. The overhead is a C4/C5 stream and the highly aromatic C6+ stream is the bottoms product. The unconverted C4 and C5 products are typically recycled. The cracking process can handle both iso and normal olefins with equivalent efficiency. There is no need to remove isobutylene for example from the feed to maximize propylene production.\nThe isomerization and metathesis of olefins in the above-described processes currently utilizes two catalysts. For example, WO3/SiO2 catalyst may perform the metathesis reaction and MgO catalysts may perform the olefin isomerization reaction. The MgO catalyst also serves a reactive guard bed, which prevents the adsorption of poisons onto the WO3/SiO2 (metathesis) catalyst, and, therefore, prolongs its cycle length.\nThe two-catalyst system typically consists of a physical mixture of two catalysts. One is a tabletted MgO and the other is a silica supported WO3 in the form of granules. The MgO catalyst is present to act both as a reactive guard bed and to provide isomerization of 1-butene to 2-butene for the reaction of 2-butene with ethylene to ultimately form two propylene molecules. These catalysts, when loaded as a (MgO: WO/SiO2) bed with a part upper MgO bed, or an overall (MgO: WO/SiO2) mixed bed, present several issues. For example, the MgO is subject to significant sintering and loss of surface area over catalyst regenerations; the loss of surface area results in loss of MgO activity for both isomerization and as an adsorbent of poisons. As another example, the mixed co-catalyst system may be subject to uneven mixing, bed non-uniformity, flow channeling, and uneven temperature distribution in the bed, which also further affects catalyst performance. Additionally, the reactive effectiveness of a physical mixture on a kinetic basis is lower than a single catalyst where the functionality is in close proximity.\nAs described above, there is considerable interest in the processing of C4, C5, and heavier olefin streams to produce lighter olefins, such as propylene. Accordingly, there exists a significant need for catalysts and processes that may result in the efficient isomerization and metathesis of linear butenes for the production of high purity propylene from such olefin-containing streams at low cost and low energy."}
{"text": "In a wireless cellular communication system, one or multiple Downlink (DL) Common Reference Signals (CRSs) may be used for channel measurements, or for coherent demodulation and channel measurements, for a mobile terminal in a given cell. A mobile terminal is also denoted as a User Equipment (UE) in some wireless communication systems. Each CRS defines a so-called antenna port in a given cell and a common way to implement antenna ports is to associate an antenna port with a physical transmit antenna.\nThe Reference Signals (RSs) of different antenna ports should be orthogonal to each other to allow interference-free identification of corresponding propagation channel coefficients, i.e. the propagation channel from each transmit antenna to each receive antenna. The RSs are usually cell-specific to minimize interference between RSs in different cells. Without loss of generality, antenna ports are defined by CRS or cell-specific CRS throughout this document. Cell specific CRS implies that they are used by multiple UEs in a cell to measure the channel from each antenna port. Antenna ports may also be user specific, which means that they are used for measurements and/or demodulation by a specific single UE.\nThe CRSs are transmitted on exclusively reserved resources of a cell, such as time and frequency Resource Elements (RE), codes, etc. Data is not transmitted on these reserved resources to avoid interference with RSs, which would hamper the estimation of the channel propagation coefficients from that antenna port.\nIn order to be able to use DL CRSs properly and to perform standard communication with a base station, such as a eNB, the number of antenna ports used for DL channel measurements and/or DL transmission is very important information that a UE needs to know. After a UE obtain information about the number of antenna ports used in a cell, the UE will know which transmission mode is used for each physical channel, and which resources that are used for data transmission and which that are used for DL CRS. This is important information to avoid that received data is punctured by CRSs, since if a UE is not aware of all CRSs in radio resources, it will assume reception of data on those resources where there actually is a CRS transmission, and this will degrade the performance of data reception due to the interference from the CRSs.\nFurthermore, knowing the number of CRSs is also important to measure multiple channels using CRSs and to detect physical channels, etc. In the Long Term Evolution (LTE) Release-8 (Rel-8) standard, the information about number of antenna ports is embedded in the signal transmitted on a Physical Broadcast Channel (PBCH). After successful cell search procedure, the UE will obtain time and frequency synchronization with a cell, as well as the cell Identity (ID) of the cell; and then the UE begins to detect the PBCH to obtain cell-specific information and the number of antenna ports.\nIn the LTE Rel-8 standard, three types of cell-specific CRSs are supported defining: one, two and four antenna ports (3GPP TS 36.211 v8.4.0). The number of antenna ports in a cell decides the maximum number of Multiple Input Multiple Output (MIMO) transmission layers supported by a eNB in said cell. For instance, if there are four antenna ports in a cell, up to four MIMO layers transmission can be supported by the eNB. The information about the number of antenna ports is embedded into the signal transmitted on the PBCH by using different Cyclic Redundancy Check (CRC) masks to indicate the number of antenna ports. As a UE has no prior information about the number of used antenna ports in a cell, i.e. the used CRC mask of the transport block of the PBCH, the UE has to make blind detection of that information, which means that it has to check all possible CRC masks and select the mask that is the most probable conditioned on the received PBCH signal.\nThe operation of embedding information about the number of antenna ports into PBCH at transmitter and the corresponding blind detection of PBCH at receiver for a LTE Rel-8 system will be described in the following.\nFirstly, at the transmitter, the entire transport block bits of PBCH a0, a1, . . . aA-1 is used to calculate the CRC parity bits p0, p1, . . . pL-1, where A is the size of the transport block (i.e. the number of information bits) and L is the number of CRC parity bits which is set to 16 in the LTE Rel-8 standard. Secondly, according to the antenna port configuration of the cell, the CRC parity bits are scrambled by a sequence having length 16, x0n, x1n, . . . x15n corresponding to a number of antenna ports n, where n=1, 2 or 4. After scrambling, the masked CRC parity bits will be c0, c1, . . . c15, where ci=(pi+xin) mod 2, i=0, 1, . . . , 15. The mapping relation between the three scrambling sequences and the number of antenna ports according to the LTE standard is shown in Table 1 below.\nTABLE 1CRC mask sequences for PBCH in LTE Rel-8Number ofantenna portsCRC mask sequences1<0, 0, 0, 0, 0, 0, 0, 0, 0, 0, 0, 0, 0, 0, 0, 0>2<1, 1, 1, 1, 1, 1, 1, 1, 1, 1, 1, 1, 1, 1, 1, 1>4<0, 1, 0, 1, 0, 1, 0, 1, 0, 1, 0, 1, 0, 1, 0, 1>\nThen, the masked CRC parity bits are attached to the tail of the transport block bits of the PBCH to obtain the information bits to be transmitted as a0, a1, . . . , aA-1, c0, c1, . . . c15.\nFinally, a set of operations including channel coding, rate matching, modulation and resources mapping are performed on the information bits a0, a1, . . . , aA-1, c0, c1, . . . c15.\nIn the case of one antenna port, the modulation symbols are directly mapped to the reserved resources on antenna port 0; in the case of two antenna ports, a transmit diversity scheme known as Space Frequency Block Coding (SFBC) is performed on the modulation symbols, and the output of the SFBC is mapped to the reserved resources on antenna port 0 and 1, respectively; in the case of four antenna ports, SFBC combined with Frequency Switching Transmit Diversity (FSTD) is performed on the modulation symbols, and the output of SFBC+FSTD is mapped to the reserved resources on antenna port 0, 1, 2 and 3, respectively. It should be observed from the above description that the information about the number of antenna ports is implicitly embedded into the PBCH.\nAt the receiver, the corresponding inverse operations to find the number of antenna ports are done by a UE who is accessing to the cell. As the UE knows that there are three hypothesises possible regarding the number of antenna ports (i.e. one, two or four antenna ports) the UE performs blind detection of the PBCH, which is illustrated in FIG. 1. In the procedure of blind detection, SFBC or SFBC+FSTD decoding, demodulation, channel decoding and CRC detection are all standard operations, so the details of them will not be further described, but the operation of removing the CRC mask will be explained in the following disclosure.\nAssuming that the output of channel decoding is â0, â1, . . . , âA-1, ĉ0, ĉ1, . . . ĉ15, where the last 16 bits of information ĉ0, ĉ1, . . . ĉ15 are the CRC parity bits scrambled with a CRC mask corresponding to information about the number of antenna ports, as mentioned above. When blind detection of the PBCH is performed, the CRC parity bits are de-scrambled with an assumed CRC mask (i.e. removing the CRC mask) in the following way:{tilde over (c)}i=(ĉi+xin)mod 2where i=0, 1, . . . , 15;                n is the number of antenna ports,        xin is the defined CRC mask corresponding to n antenna ports        \nIf the assumed CRC mask is the same as the actual CRC mask used at transmitter, the above operation will completely remove the CRC mask embedded into CRC parity bits, and the probability of correct detection is increased.\nAs mentioned, in the LTE Rel-8 system up to four antenna ports can be supported on the DL. The LTE-Advanced (LTE-A) system of Release 10 (Rel-10) and beyond are supposed to be an extension of LTE system in which up to eight layers transmission (possibly even more layers for releases beyond Rel-10) will be supported to further increase system performance, such as peak data rate, cell average spectrum efficiency, etc (3GPP TR 36.814 v1.0.0). In order to support up to eight layers transmission more antenna ports than the antenna ports supported in LTE Rel-8 must be defined in a LTE-A communication system.\nIn addition, to fulfil LTE-A backwards compatibility requirement, it should still be possible for a LTE-A cell to also serve LTE UEs. In order to enable LTE UEs to operate in a LTE-A system, the antenna ports defined in LTE should also be supported in a LTE-A system, i.e. n number of LTE CRSs should also exist in a LTE-A system, where n=1, 2 or 4; and LTE UEs use LTE CRSs for coherent demodulation and channel measurement as in the LTE system, while LTE-A UEs may also use these LTE CRS for demodulation of the control channels, such as PBCH and Physical Downlink Control Channel (PDCCH).\nHence, in a LTE-A system there will be a number of LTE antenna ports used for transmitting LTE data/control and/or LTE-A control information; in addition, it is also possible to define a number of additional antenna ports used only for supporting LTE-A data transmission. For all antenna ports in a LTE-A system, the new defined additional antenna ports are denoted as LTE-A antenna ports, and the number of LTE-A antenna ports could be zero. Therefore, a communication system which can serve both LTE UEs and LTE-A UEs is needed. Also, since the number of LTE and LTE-A antenna ports may be different in a LTE-A cell, a question is how to signal the existence of LTE-A antenna ports to LTE-A UEs in way that is transparent to the reception of the number of LTE antenna ports.\nIt is thus clear that a LTE-A eNB may need to enable additional CRSs (antenna ports) for measurements and/or demodulation compared to the antenna ports defined by LTE CRSs (i.e. n=1, 2 or 4). Up to eight additional antenna ports (CRSs) may be needed for LTE-A UEs in Rel-10. These additional CRSs are denoted Channel State Information-RSs (CSI-RSs). It is thus another question how to signal the number of LTE-A antenna ports or the number of CSI-RS to LTE-A UEs in way that is transparent to LTE UEs.\nTherefore, a signalling method which is backwards compatible to enable LTE UEs to obtain the number of LTE antenna ports, while the signalling of LTE-A antenna ports should be transparent to LTE UEs is needed in the art."}
{"text": "Field of the Invention\nThe invention relates to a semiconductor device, and particularly relates to a transistor structure.\nDescription of Related Art\nThe high voltage transistor device (such as lateral diffused metal-oxide semiconductor (LDMOS) etc.) is broadly applied to various integrated circuits (ICs). However, the charge and the electric field is easy to concentrate at the sharp edge or the sharp point of the high voltage transistor device, and thus the leakage current is occurred and the electrical performance is reduced."}
{"text": "In a typical personal computer, there exists a processor, a basic input-output system (BIOS), a main memory, and a hard drive that stores an operating system and one or more application programs. The BIOS usually takes the form of executable instructions stored on a read-only memory (ROM). During a regular boot-up process, the processor accesses and executes the instructions in the BIOS, and under direction of the BIOS, the processor implements the booting process. The BIOS causes the processor to perform some low-level setup functions to prepare the computer for regular operation. After the setup functions are performed, the BIOS causes the processor to load and execute the operating system stored on the hard drive. By doing so, the BIOS in effect transfers control from itself to the operating system. After the operating system is loaded and executed by the processor, the boot-up process is complete and the computer is ready for operation.\nThe ROM (read-only-memory) on which the BIOS resides (referred to hereinafter as the BIOS ROM) is typically quite small in size. Despite this small size, however, it has been observed that one or more auxiliary programs may be stored on the BIOS ROM. These auxiliary programs may be executed during the booting process prior to and even in lieu of the operating system on the hard drive to provide certain desired functionalities. An example of an auxiliary program is the self-contained browser disclosed in U.S. patent application Ser. No. 09/449,065 entitled “Self Contained Browser Architecture” filed on Nov. 24, 1999, which is now abandoned, the contents of which are incorporated herein by this reference. Auxiliary programs are generally self-contained, meaning that they comprise all of the components that they need to operate. As a result, they do not need the operating system on the hard drive to function. This can be quite advantageous because even if the hard drive fails or the operating system becomes corrupted, the auxiliary programs are not affected. They can still function. Hence, the auxiliary programs are impervious to many system failures.\nCurrently, auxiliary programs are executed, if at all, during the boot-up process prior to the transfer of control to the operating system. A mechanism for invoking an auxiliary program during a boot-up process is disclosed in U.S. Pat. No. 7,082,526 entitled “Mechanism for Intuitively Invoking One or More Auxiliary Programs During A Computer Booting Process” filed on Mar. 14, 2003, the contents of which are incorporated herein by this reference. Once the operating system assumes control, however, the auxiliary programs typically can no longer be accessed. This means that if a user wishes to execute an auxiliary program after the operating system has already loaded and executed, the user has to reboot the computer. This can be an inconvenient and time-consuming process."}
{"text": "Opto-electronic devices that make use of organic materials are becoming increasingly desirable for a number of reasons. Many of the materials used to make such devices are relatively inexpensive, so organic opto-electronic devices have the potential for cost advantages over inorganic devices. In addition, the inherent properties of organic materials, such as their flexibility, may make them well suited for particular applications such as fabrication on a flexible substrate. Examples of organic opto-electronic devices include organic light emitting devices (OLEDs), organic phototransistors, organic photovoltaic cells, and organic photodetectors.\nMore details on OLEDs can be found in U.S. Pat. No. 7,279,704, which is incorporated herein by reference in its entirety.\nVarious ways to deposit the organic materials used to fabricate organic devices are known, such as vacuum thermal evaporation, solution processing, organic vapor phase deposition, and organic vapor jet printing."}
{"text": "Airbags may be mounted within a vehicle and deployed so as to prevent a vehicle occupant from impact with a vehicular structure during a collision event. Some airbags suffer from one or more drawbacks or may perform less than optimally in one or more respects. Certain embodiments disclosed herein can address one or more of these issues."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a printer, and more particularly to a printer which makes a print by scanning a sheet of photo-sensitive material with a beam of light emitted from a light source of light emitting diodes, or the like.\n2. Description of Related Art\nThere is a conventional printer which makes a print by scanning a sheet on which a photo-sensitive material is applied (hereinafter referred to as a photo-sensitive sheet) with a beam of light emitted from a light source of light emitting diodes (LED), or the like. In the conventional printer, for main scanning, a light beam from the LED light source is deflected by a scanning mirror which is oscillated by a cam mechanism so that the light beam can traverse the photo-sensitive sheet in the form of a line of light (hereinafter referred to as a line beam) in a main scanning direction. For sub-scanning, a unit (a stage) on which the photo-sensitive sheet is set (hereinafter referred to as a sub-scanning unit) is horizontally moved in a sub-scanning direction which is perpendicular to the main scanning direction so that the photo-sensitive sheet can be scanned with the line beam from the LED light source in the sub-scanning direction. Thus, the photo-sensitive sheet can be scanned with the light beam from the LED light source in both the main scanning direction and the sub-scanning direction, and thereby an image is printed on the photo-sensitive sheet.\nIn the conventional printer, a face cam is used in the cam mechanism for oscillating the scanning mirror, and a cam follower is pressed against a cam face of the face cam by a biasing spring, so that the cam follower can follow the cam face of the face cam, which is rotated by a motor. This cam mechanism, however, has a problem in that a contact friction between the face cam and the cam follower and a load torque of the drive motor are great. Moreover, the load torque fluctuates according to the movement of the cam follower which moves up and down on the cam face of the face cam, and the fluctuation of the load torque has bad effects on the print. If the conventional printer using the cam mechanism makes a print at higher speed, the force of the biasing spring to press the cam follower must be increased in proportion to the square of the rate of revolution of the face cam, and the driving force of the diving motor must be increased in proportion to the cube of the rate of revolution of the face cam.\nIn the conventional printer, the scanning mirror is oscillated by the cam mechanism shown in FIG. 19, and a light beam in the form of a line in the main scanning direction is generated. As shown in FIG. 19, a swing pin 304 is attached to the end of a swing lever 302, and the swing pin 304 is biased against a face cam 300, which is rotated by a motor (not shown), by a biasing spring (not shown). A scanning mirror 306 is fixed to the swing lever 302. A light beam enters from the right side in FIG. 19, and is reflected downward by the scanning mirror 306. When the motor rotates the face cam 300, the swing lever 302 swings up and down, and the light beam changes its reflecting directions to the right and left to thereby generate a line beam.\nIn the conventional printer, the line beam in the main scanning direction falls on the photo-sensitive material while the sub-scanning unit, on which the photo-sensitive material is set, is moving in the sub-scanning direction as described above. Thus, the line beam which is reflected by the mirror 306 as shown in FIG. 20(A) describes forward (left to right) scanning lines and backward (right to left) scanning lines on the photo-sensitive sheet as shown in FIG. 20(B), and both scanning lines are inclined because the photo-sensitive sheet moves (upward in the drawing). For this reason, the forward and backward scanning lines cannot be parallel to one another, and the resolution of the print can be deteriorated. Hence, the conventional printer makes a print using only the forward scanning lines, and then the number of scanning lines is reduced by half. In this case, the quality of the print is lowered, and the quantity of light is poor. To solve these problems, there is a printer which uses a fast-forwarding mechanism to oscillate the scanning mirror, but there is a disadvantage that the printer shakes and makes noise.\nIn the conventional printer, if one motor drives the cam mechanism which oscillates the scanning mirror and the sub-scanning unit on which the photo-sensitive material is set, the cam mechanism oscillates the scanning mirror when the sub-scanning unit returns to the start printing position at a high speed after moving to make a print (for example, because the swing pin, which swings the scanning mirror, abuts on a cam member such as the face cam). Thus, a large load is applied to the motor, and the cam mechanism makes loud noise while rotating at a high speed."}
{"text": "In traditional metal insulator semiconductor (MIS) dynamic random access memory (DRAM) devices, high temperature nitridation before the high dielectric constant (k) deposition is needed to prevent bottom electrode oxidation during the high-k deposition. It is noted that a high dielectric constant (k) is above about 3.9.\nThe major concern is to eliminate native oxide and to increase capacitance. However, a high temperature anneal is not acceptable for embedded DRAM because logic performance will be degraded.\nU.S. Pat. No. 6,580,115 B2 to Agarwal describes a capacitor electrode for integrating high-k materials (wherein high-k materials have a dielectric of greater than about 20).\nU.S. Pat. No. 5,663,098 to Creighton et al. describes a method for deposition of a conductor in integrated circuits.\nU.S. Pat. No. 4,751,101 to Joshi describes low stress tungsten films by silicon reduction of WF6."}
{"text": "An ATM (Automatic Teller Machine) has been widely used in various commercial banks and for postal saving. The ATM operates in the 24-hour self-help manner, which brings people great convenience.\nIn an ATM automatic banknote adding system, all of cash recycling ATMs, cash deposit ATMs and banknote adding machines share a same kind of banknote boxes. The banknote box itself has an information memory for storing information on banknotes in the banknote box. Each banknote box has a physical ID number which is globally unique, so as to be able to distinguish from each other. The automatic banknote adding system may directly transmit a banknote box which is full of banknotes in a cash deposit ATM to a cash outputting ATM, and the cash outputting ATM automatically gets the information on the banknote in the banknote box, so as to automatically add banknotes. The banknote box may be changed at any moment at will between all banknote processing machines in the automatic banknote adding system, and information on each replacing banknote box is uploaded to a server of a bank through each ATM so as to record corresponding data, thus no man-made count is needed, so that the efficiency of processing banknote is greatly improved.\nPresently, if an ATM is used for cash withdrawal, the ATM may send unqualified banknotes to a recovery part. However, information (including the number and face values of the banknotes) on the banknotes sent to the recovery part cannot be correctly determined. Therefore, it is impossible to accurately obtain the information on the residual banknotes in the banknote box.\nFor example, if a certain banknote is in an unqualified status such as an overlapped banknote or an adhesion banknote) when the banknote is output from a cash recycling ATM, it can be determined that the banknote is in the abnormal status by measuring the thickness or width, and thus the banknote is recovered. However, it is very difficult to determine the number of the recovered banknotes (taking the overlapped banknotes as an example, it is impossible to determine whether three or two banknotes are overlapped; and if there is adhesive substance on one banknote, the banknote may also be mistakenly determined as two banknotes). Therefore, it is unable to accurately track the information on the residual banknotes in the banknote box. To accurately obtain the information on banknotes in the banknote box is an important prerequisite for the ATM to automatically add banknotes."}
{"text": "Various proposals have been offered regarding the protection of vehicles against theft Japanese Unexamined Patent Publication No. Hei 8(1996)-91176, for example, teaches an arrangement in which engine control means is provided that checks a key ID code against a reference ID and enables engine starting when the two are in a relationship defined beforehand, and the engine control means is connected with various drive/control means to be controlled by connectors having a unique shapes in accordance with the configuration thereof, whereby the connectors of an engine control means of another configuration and the connectors of the various drive/control means cannot be properly electrically connected.\nHowever, this proposed system does not always function as an effective immobilizer because the engine can be easily started by obtaining another genuine engine control means whose ID code is known and replacing the engine control means installed in the vehicle.\nAn object of the present invention is therefore to provide a vehicle electronic control system with anti-theft capability and a method of controlling a vehicle engine which overcome this problem of the prior art.\nAnother object of the present invention is to provide a method of controlling a vehicle engine capable of preventing vehicle theft."}
{"text": "The present invention relates to a process for the preparation of a low viscosity alkyl toluene or alkyl xylene sulfonate.\nAlkyl toluene and alkyl xylene sulfonates have utility as surfactants for use in enhanced oil recovery. However, these sulfonates are often highly viscous and therefore can be difficult to handle and transport except in highly dilute form. Since the selection of a particular surfactant which can be utilized in petroleum recovery techniques is based in part upon economic considerations, it would be advantageous to have a method for obtaining a low viscosity alkyl toluene or alkyl xylene sulfonate.\nU.S. Pat. No. 3,957,671 describes a low viscosity detergent acid mix containing alkyl benzene sulfonic acid which is prepared by sulfonating alkyl benzene in the presence of benzoic acid."}
{"text": "1. Field of the Invention\nThe present invention relates to expanding markets for products and services by making resource placement decisions in an objective manner, and deals more particularly with techniques for comparing one or more locations to a set of criteria that are directed toward identifying market strengths and weaknesses of the location(s) as part of an overall value chain.\n2. Description of the Related Art\nMany countries or localities represent an untapped or an emerging market for the sale and adoption of information technology (“IT”). Often, these same countries are simultaneously experiencing a high growth rate in the availability of skills supportive of IT. Connections between a company's business and marketing strategy (and its related decisions) and its resource placement strategy (and its related decisions) may be revealed by examining and understanding the value chain in which the company participates.\nA value chain is defined as a sequence or network of transactions and mutually-beneficial relationships occurring between companies in the delivery of value to an end customer. Value chains have also been referred to as supply chains and value nets. Economists as far back as Adam Smith have described in great detail how value chains work, especially with regard to how such chains allow for and encourage specialization of business roles, and how they also drive more responsive business systems over time.\nAs an example, consider a simple value chain that corresponds to a hardware supply chain, wherein a wholesale hardware manufacturer supplies a hardware product directly to a retailer who in turn sells it to an end customer. The hardware manufacturer receives value from the retailer because the retailer enables the manufacturer's products to reach the marketplace; this link in the value chain provided by the retailer supports and sustains further orders for the products provided by the manufacturer, thereby causing an increase not only in the manufacturer's revenue but also in the revenue of the retailer. The retailer receives value from the hardware manufacturer in terms of obtaining products that will meet needs of customers in this marketplace. The end customer receives value from the retailer, who provides the customer with an opportunity to buy the needed product, and in turn, the retailer receives value from the customer from the sale of the products. Other benefits may also result, such as increased market presence or brand awareness, and so forth. Symbiotic relationships thus exist between entities in the value chain.\nMore complicated value chains, which may be recursive in some cases, are seen in the IT field. FIG. 1 illustrates representative entities in a sample IT-related value chain. A plurality of components 110 may be required in this particular chain, for example, and these components may be adapted for use with a plurality of operating systems 120. In the more general case, a number of entities (see, for example, reference number 130 in FIG. 1) may be involved in creating one or more IT solutions 140, which may be marketed through one or more distribution channels 150.\nWith these simple examples, it should be clear how companies become part of a value chain and also how they become dependent on both the upstream and downstream vitality and growth of other links in the chain. No company in a mature and free market is likely to own the entire chain of value. In fact, economists use the complexity of a value chain as one indicator of a market's (and a country's) economic maturity. It therefore follows that a company's ability to expand its market presence by delivering value to emerging and immature markets is influenced and limited by its ability to develop a rich and differentiated value net within these emerging markets.\nAccordingly, it is desirable to provide techniques that leverage a value chain analysis approach for determining optimal resource placement, particularly in emerging IT markets."}
{"text": "1. Field\nThe following description relates to an equal model conservation technique for a pipeline processor.\n2. Description of the Related Art\nAn instruction is processed at several stages in a pipeline processor. For example, an instruction may be processed at a fetch stage, a decoding stage, an execution stage, and a write-back stage. In the pipeline processor, because a plurality of instructions sequentially pass through the individual stages of the pipeline while overlapping each during execution, programs can be processed at high speeds.\nExamples of a pipeline processor include an equal model processor and a less-than-or-equal (LTE) processor. In the pipeline processor, a maximum latency may be set for each instruction. For example, a value of maximum latency for each instruction may depend on a processor type.\n“Maximum latency” indicates a maximum time period for which an operation corresponding to each instruction is executed and the result of the execution is written in a target register. In the equal model processor, the result of the execution of an instruction is written in the target register after the set maximum latency elapses. Meanwhile, in the LTE model processor, the result of execution of an instruction may be written in the target register before set maximum latency elapses. Accordingly, the equal model processor allows flexible scheduling through multiple assignment of a target register because the result of the execution of an instruction is not written in the target register until maximum latency elapses. Alternatively, in the LTE model processor, because the result of execution of an instruction can be written in a target register before maximum latency elapses, multiple assignment of the target register is limited."}
{"text": "The invention relates to a runflat device for a motor vehicle, enabling a considerable distance to be traveled at relatively high speed with a tuneless tire that is partially or totally deflated.\nIn general, for fitting to civilian vehicles, runflat devices need to be suitable for being mounted on standard wheel rims, which means one-piece rims of the drop-center type.\nPresently-known runflat devices are generally constituted by a rolling ring which is clamped around the rim of the wheel inside the tire. The ring is constituted either by a single, relatively flexible piece from which a splice has been removed, or else by at least two rigid pieces in the form of circular arcs or sectors. To mount the ring in a clamped configuration around the rim, it is necessary to provide assembly and clamping fastener means between the two free ends of the split ring or between the facing ends of the ring sectors. Assembly and clamping fastener means are rigid and they are constituted by mechanical elements such as nuts-and-bolts, for example.\nUnfortunately, experience shows that it is such rigid fastener means that constitute the weak links in runflat devices (e.g. Due to fatigue phenomena). Furthermore, mounting a runflat device is an operation that requires the operator to be trained, in particular because of the small amount of space available for the tool that is required for clamping the ring onto the rim.\nAccording to document FR-97/13618 in the name of the Applicant, the runflat device is made up in such a manner as to eliminate assembly and clamping fastener means. It comprises at least one open inner ring of relatively rigid material for mounting on the rim, and a continuous outer ring that is substantially in extensible and that is engaged on the inner ring so as to clamp it and fix it to the rim.\nIn an improvement disclosed in document FR-98/04225, likewise in the name of the Applicant, the ring is made up of three pieces, namely: two first annular pieces that are open at respective slices and that are designed to be mounted one on the other around the rim, and a third annular piece that is continuous and coaxial about the two preceding pieces, said third piece serving to clamp the system built up in this way onto the rim.\nAn object of the invention is specifically to avoid problems associated with assembly and clamping fastener means between the two free ends of a ring or between sectors of a ring.\nTo this end, the invention provides a runflat device for a motor vehicle, the device comprising a ring for mounting inside a tire on a one-piece non-standard wheel rim that presents a drop center, an inner flange and an outer flange, said ring being made as at least two sectors each having a radially inner zone and a radially outer zone, wherein the sectors of the ring are independent of one another, and wherein each sector is fixed to the rim via its radially inner zone.\nIn general, the radially inner zone of each sector of the ring is terminated by a circularly actuate footing designed to bear against the bottom of the drop center, and said footing presents an inner lateral flange and an outer lateral flange, the inner flange being received in a lateral groove in the drop center while the outer flange matches the profile of the drop center and cooperates with the means for fixing the sector of the ring on the rim.\nIn the invention, each sector of the ring is secured to the wheel rim by means of screws or by means of a clamping belt.\nIn a first embodiment of the invention, fixing screws are used which pass through the outer flange of the footing of each sector of the ring and penetrate into tapped blind holes formed in the thickness of the wheel rim from the bottom of the drop center.\nAdvantageously, the tapped blind holes are machined obliquely in the thickness of the rim to facilitate screw-fitting operations.\nIn a second embodiment of the invention, fixing screws are used which pass through tapped holes machined laterally through the rim and opening out into the drop center so as to bear in clamped manner against the outside faces of the sectors of the ring.\nIn addition, the inner flanges of the footings of the sectors of the ring bear against an O-ring placed in the lateral groove of the drop center and designed to accommodate the lateral tolerances of the various parts of the assembly, while each fixing screw has an enlarged head which bears against the rim via a sealing ring.\nIn a third embodiment of the invention, the means for fixing the sectors of the ring on the rim are constituted by a clamping belt which surrounds the outer flanges of the sectors of the ring."}
{"text": "1. Field of the Invention\nThis invention relates generally to processor-based systems, and, more particularly, to error recovery in a directory cache of a distributed, shared-memory processor-based system.\n2. Description of the Related Art\nBusinesses typically rely on network computing to maintain a competitive advantage over other businesses. As such, developers, when designing processor-based systems for use in network-centric environments, may take several factors into consideration to meet the expectation of the customers, factors such as functionality, reliability, scalability, and performance of such systems.\nOne example of a processor-based system used in a network-centric environment is a mid-range server system. A single mid-range server system may have a plurality of system boards that may, for example, be configured as one or more domains, where a domain, for example, may act as a separate machine by running its own instance of an operating system to perform one or more of the configured tasks.\nA mid-range server, in one embodiment, may employ a distributed shared memory system, where processors from one system board can access memory contents from another system board. The union of all of the memories on the system boards of the mid-range server comprises a distributed shared memory (DSM).\nOne method of accessing data from other system boards within a system is to broadcast a memory request on a common bus. For example, if a requesting system board desires to access information stored in a memory line residing in a memory of another system board, the requesting system board typically broadcasts on the common bus its memory access request. All of the system boards in the system may receive the same request, and the system board whose memory address ranges match the memory address provided in the memory access request may then respond.\nThe broadcast approach for accessing contents of memories in other system boards may work adequately when a relatively small number of system boards are present in a system. However, such an approach may be unsuitable as the number of system boards grows. As the number of system boards grows, so does the number of memory access requests, thus to handle this increased traffic, larger and faster buses may be needed to allow the memory accesses to complete in a timely manner. Operating a large bus at high speeds may be problematic because of electrical concerns, in part, due to high capacitance, inductance, and the like. Furthermore, a larger number of boards within a system may require extra broadcasts, which could further add undesirable delays and may require additional processing power to handle the extra broadcasts.\nDesigners have proposed the use of directory caches in a distributed shared memory systems to reduce the need for globally broadcasting memory requests. Typically, each system board serves as home board for memory lines within a selected memory address range, and where each system board is aware of the memory address ranges belonging to the other system boards within the system. Each home board generally maintains its own directory cache for memory lines that fall within its address range. Thus, when a requesting board desires to access memory contents from another board, instead of generally broadcasting the memory request in the system, the request is transmitted to the appropriate home board. The home board may consult its directory cache and determine which system board is capable of responding to the memory request.\nDirectory caches are generally effective in reducing the need for globally broadcasting memory requests during memory accesses. However, as would be expected, the effectiveness of the directory caches depends in part on the directory caches being properly operational while the system is running. An inoperable or a partially inoperable (i.e., functioning but with one or more errors) directory cache may sometimes go undetected for extended periods of time, and may thereby adversely affect the overall operation of the system."}
{"text": "There is known a semiconductor package in which a semiconductor chip (semiconductor element) is covered (encapsulated) with a resin-made encapsulating (sealing) material. The encapsulating material for the semiconductor chip is produced by molding a resin composition through, e.g., a transfer molding method.\nA process for producing the resin composition from a powder material thereof includes a surface treatment process in which a processing liquid such as a coupling agent is allowed to adhere to surfaces of inorganic particles such as silica particles (filler material) contained in the powder material of the resin composition.\nIn this surface treatment process, the silica particles are injected into a chamber from an inlet port thereof, and the coupling agent is also sprayed into the chamber from an opposite side with respect to the inlet port while the silica particles are spirally swirled in the chamber (for example, see patent document 1). In this configuration, a part of the coupling agent is contacted with the silica particles upon spraying the coupling agent and the remaining coupling agent is swirled and contacted with the other silica particles in the chamber during swirling of the silica particles so that the coupling agent adheres to the surfaces of the silica particles. As a result, a mixing property of resin particles constituting the resin composition and the silica particles is enhanced, thereby easily mixing the silica particles and the resin particles since the silica particles are coated with the coupling agent.\nHowever, according to the above method, in the surface treatment process, there is a problem in that the silica particles are aggregated. In this case, it is difficult to uniformly mix the silica particles (inorganic particles) with the resin particles (resin composition).\nPatent Document 1: JP-A 2003-275555.\nIt is an object of the present invention to provide a particle production apparatus, a particle production method and a method for producing a semiconductor encapsulating resin composition, by which a processing liquid is allowed to adhere to inorganic particles contained in a powder material reliably while aggregation of the inorganic particles is suppressed.\nIn order to achieve the object, one aspect of the present invention is directed to a particle production apparatus which is configured so that a processing liquid adheres to a surface of each of inorganic particles contained in a powder material, the particle production apparatus comprising:\na processing section in which the processing liquid is allowed to adhere to the surface of each of the inorganic particles;\na chamber connected to the processing section at a downstream side thereof in which the powder material is separated from gas carrying the powder material;\na powder material supply device which supplies the powder material into the processing section; and\na processing liquid spraying device which sprays the processing liquid as droplets to the powder material just after the powder material being supplied into the processing section,\nwherein after the processing liquid sprayed to the powder material adheres to the surface of each of the inorganic particles contained in the powder material in the processing section, the powder material is transferred into the chamber so that the powder material is separated from the gas.\nIn the particle production apparatus of the present invention, it is preferred that the processing liquid spraying device comprises: a nozzle; a gas supply device for supplying the gas having a pressure being equal to or higher than 0.3 MPa to the nozzle; and a processing liquid supply device for supplying the processing liquid into the nozzle, and\nwherein the processing liquid is sprayed from the nozzle into the processing section due to the pressure of the gas.\nIn the particle production apparatus of the present invention, it is preferred that the processing liquid supply device is a pump.\nIn the particle production apparatus of the present invention, it is preferred that the droplets of the processing liquid sprayed from the processing liquid spraying device include small-sized droplets having a particle size being equal to or less than 20 and the processing liquid spraying device is configured so that a ratio of the small-sized droplets to the total of the sprayed droplets becomes equal to or more than 80 wt %.\nIn the particle production apparatus of the present invention, it is preferred that the particle production apparatus is configured so that a swirl stream of the gas is generated in the chamber by spraying the processing liquid using the processing liquid spraying device.\nIn the particle production apparatus of the present invention, it is preferred that the chamber has a cylindrical shape and the processing section has a central axis and an outlet port connected to the chamber,\nwherein the processing section is connected to a side portion of the chamber such that the central axis of the processing section is inclined with respect to a direction of a radius of the chamber which passes through the outlet port of the processing section, and\nwherein the particle production apparatus is configured so that a swirl stream of the gas is generated in the chamber by spraying the processing liquid using the processing liquid spraying device.\nIn the particle production apparatus of the present invention, it is preferred that the inorganic particles are constituted of silica and the processing liquid is a coupling agent.\nIn order to achieve the object, another aspect of the present invention is directed to a particle production method by which a processing liquid is allowed to adhere to a surface of each of inorganic particles contained in a powder material, wherein the powder material is supplied into a processing section, and the processing liquid is sprayed as droplets to the powder material just after the powder material being supplied into the processing section, thereby allowing the processing liquid to adhere to the surface of each of the inorganic particles contained in the powder material, and then the powder material is transferred into a chamber so that the powder material is separated from gas carrying the powder material.\nIn the particle production method of the present invention, it is preferred that the processing liquid and the gas having a pressure being equal to or higher than 0.3 MPa are supplied into a nozzle so that the processing liquid is sprayed from the nozzle into the processing section.\nIn the particle production method of the present invention, it is preferred that an average particle size of the inorganic particles is in the range of 0.5 to 100 μm.\nIn the particle production method of the present invention, it is preferred that when the processing liquid is sprayed to the powder material so that the processing liquid adheres to the surface of each of the inorganic particles contained in the powder material, there is a case that the powder material includes aggregated particles each comprised of the inorganic particle and the adhering processing liquid wherein each aggregated particle having a particle size of 150 μm or more, but a ratio of such aggregated particles to the whole of the powder material is equal to or less than 1 wt % of the powder material.\nIn order to achieve the object, the other aspect of the present invention is directed to a method for producing a semiconductor encapsulating resin composition comprising:\nallowing the processing liquid to adhere to the surface of each of the inorganic particles contained in the powder material by the particle production method of the present invention; and\nmixing a composition containing a powder material constituted of a resin and the powder material containing the inorganic particles to which the processing liquid adheres."}
{"text": "In the medical field, it is sometimes the case that when a region of interest exists on an image obtained by a given modality, a corresponding region is identified on an image obtained by another modality, and diagnosis is then performed based on the comparison between them. If imaging posture differs between modalities, the shape of the object differs at the times of imaging. This makes it difficult to identify the object. Under the circumstance, attempts have been made to estimate the deformation of an object between modalities (that is, alignment between images accompanying deformation). This makes it possible to estimate the position of a corresponding region based on the position information of a region of interest and to generate an image by deforming one image so as to have the same shape as that of the other image.\nNPL 1 discloses a technique of facilitating the comparison between the shapes of an object before and after its deformation by normalizing the shape of the object accompanying deformation. More specifically, there is disclosed a method of calculating the geodesic distance matrix of the surface shape of an object and normalizing the surface shape by using a multidimensional scaling method using the calculated matrix. This method can normalize the shape of an object which deforms without changing the geodesic distance matrix of the object surface so as to allow direct comparison between the shapes of the object before and after the deformation. This makes it possible to easily compare the deformed shapes of an object accompanying deformation and recognize the object based on its shape."}
{"text": "Mobile computing devices are increasingly being adopted by mainstream users. The combination of easy portability, increasing network availability, and large local storage capabilities will result in mobile devices becoming the primary data repository for users. This use of mobile devices to carry the user's core data is a natural evolution of varied uses of previous generations of mobile devices. For example, cellular phones, media players, navigation devices, personal digital assistants (PDAs) and the like, were used to carry specialized data related to the device's primary function, e.g., contact data, music/video, maps and geographical data, notes, task lists, etc. Mobile devices have evolved that are capable of performing all of those functions on a single apparatus, and as a result the data related to those functions is also stored on the device.\nAs mobile devices have gained these various capabilities, the devices are called upon to access a wide variety of local and remote content. Local content may include any files stored on a device's persistent storage or directly accessible via a peripheral interface. Remote content may include data that is accessible via networks, including infrastructure or ad-hoc networks accessible in the home, office, and/or the Internet. As a result of the ubiquity of networking, many uses now blur the lines between local and remote content. For example, a music subscription service may offer music downloads and/or streaming, and be set up in such a way that the user need not know whether a currently playing song is locally stored or streaming over a network.\nEven though there may be convergence between local and remote content for a given user application, the way that a user searches for such information may be still be configured to search for a particular type of data in a particular domain. This can lead to confusion in cases where the user is not particularly sure in which domain the target may reside. Further, such specialized searches may be inefficient, particularly in a reduced interface mobile device. Having to type in a query in different programs and in search contexts may be tiresome with a mobile input device. Similarly, trying to view and assimilate results on a small screen may be difficult."}
{"text": "Tetrahydrofolate synthetase (to be referred to as C1-THFS hereinafter) is an enzyme which synthesizes a tetrahydrofolic acid derivative that provides C1 group necessary for various metabolic reactions (Non-patent Reference 1; hereinafter, tetrahydrofolic acid is referred to as TF, and tetrahydrofolic acid derivative as TF derivative). Illustratively, TF derivatives give C1 group to the biosynthesis reactions of purine, thymidylic acid histidine, pantothenic acid and the like. That is, TF and TF derivatives are deeply concerned in nucleic acid metabolism, amino acid metabolism and the like. Thus, TF and TF derivatives are found in tissues where cell division is vigorously carried out and essential for cell multiplication and growth.\nC1-THFS is a tri-functional enzyme which has 3 types of functions. Illustratively, C1-THFS has the functions of 10-formyl-THF synthetase (EC 6.3.4.3), 5,10-methenyl-THF cyclohydrolase (EC 3.5.4.9) and 5,10-methylene-THF dehydrogenase (EC 1.5.1.5). By exerting these functions, C1-THFS accelerates synthesis of TF derivatives which are necessary for various metabolic reactions.\nAnalysis of C1-THFS has been carried out using human, mouse, yeast and the like eucaryote, Escherichia coli and the like procaryote and the like various organism species as the object (Non-patent References 2 to 5). Regarding yeast, the presence of a C1-THFS which functions in its cytoplasm and mitochondria is known. Also, regarding human, the presence of a C1-THFS which functions in the cytoplasm of its cells is known.\nHowever, the presence of a C1-THFS which functions in the mitochondria of human cells is not known, and there is only one report of Non-patent Reference 19. In addition, a human C1-THFS gene whose expression is elevated in colon cancer tissue in comparison with normal colon tissue is not known, too.\nThe references cited in this description are listed in the following.\nNon-patent Reference 1\nHum D W et al., The Journal of Biological Chemistry, 1988, vol. 263, no. 31, pp. 15946-15950.\nNon-patent Reference 2\nStaben, C et al., The Journal of Biological Chemistry, 1984, vol. 261, pp. 4629-4637.\nNon-patent Reference 3\nShannon, K. W. et al., The Journal of Biological Chemistry, 1986, vol. 261, pp. 12266-12271.\nNon-patent Reference 4\nThigpen, A. E. et al., The Journal of Biological Chemistry, 1990, vol. 265, pp. 7907-7913.\nNon-patent Reference 5\nDev, I. K. et al., The Journal of Biological Chemistry, 1978, vol. 253, pp. 4245-4253.\nNon-patent Reference 6\nRaj S K et al., Biochemistry and Molecular Biology International, vol. 44, no. 1, pp. 89-95.\nNon-patent Reference 7\nFrohman M. A. et al., Proceedings of The National Academy of Science of The United States of America, 1988, vol. 85, no. 23, pp. 8998-9002.\nNon-patent Reference 8\nProceedings of The National Academy of Science of The United States of America, 1977, vol. 74, pp. 5463-5467.\nNon-patent Reference 9\nMethods in Enzymology, 1980, no. 65, p. 499-.\nNon-patent Reference 10\nClaros M G et al., European Journal of Biochemistry, 1996, vol. 241, no. 3, pp. 779-786.\nNon-patent Reference 11\nKim P J et al., Lancet, 2003, vol. 362, pp. 205-209.\nNon-patent Reference 12\nR. H. Giles et al., Biochimica et Biophysica Acta, 2003, vol. 1653, pp. 1-24.\nNon-patent Reference 13\nHe T C et al., Science, 1998, vol. 128, pp. 1509-1515.\nNon-patent Reference 14\nLevens D L, Genes and Development, 2003, no. 17, pp. 1071 1077.\nNon-patent Reference 15\nGarrow T A et al., The Journal of Biological Chemistry, 1993, vol. 268, pp. 11910-11916.\nNon-patent Reference 16\nNikiforov M A et al., Molecular and Cellular Biology, 2002, vol. 22, pp. 5793-5800.\nNon-patent Reference 17\nTavtigian S V et al., Molecular Biology of the Cell, 1994, vol. 5, pp. 375-388.\nNon-patent Reference 18\nVillar E, The Journal of Biological Chemistry, 1985, vol. 260, no. 4, pp. 2245-2252.\nNon-patent Reference 19\nPrasannan P et al., The Journal of Biological Chemistry, 2003, vol. 278, no. 44, pp. 43178-43187.\nNon-patent Reference 20\nSugiura et al., Biochemical and Biophysical Research Communications, 2004, vol. 315, no. 1, pp. 204-211."}
{"text": "The present invention relates to boiling heat transfer surfaces which are provided on a heat exchange wall surface to enhance nucleate boiling and thereby improve the heat transfer action of a boiling liquid such as a halocarbon refrigerant. It is well known that nucleate boiling in a heat transfer apparatus is improved when the heated wall portion through which heat is transmitted from a warm fluid to a boiling liquid has its surface made porous by the formation of cavities in the heated wall portion.\nU.S. Pat. No. 3,384,154 teaches that the radii of the cavities provided in the heated wall portion can be controlled by the size of the sintered powder coatings utilized to make the nucleate boiling surface.\nU.S. Pat. No. 3,906,604 controls the nucleate boiling cavity size by machining the metal substrate to create protrusions which are frictionally rubbed and stretched to form the cavities.\nU.S. Pat. Nos. 3,696,861 and 3,768,290 utilize bent-over tube fins to create elongated nucleate boiling cavities.\nGerman Publication No. 2,510,580 discloses heat pipe surfaces wherein the surface is first mechanically roughened by wire brushing to form close uniform scratches and then electroplated to form nucleate boiling cavities.\nU.S. Pat. No. 4,018,264 discloses a boiling heat transfer surface wherein the nucleate boiling cavities are made by electrodepositing copper dendrites on a copper surface at a high current, copper plating the dendrites or nodules and then rolling the dendritic surface to partially compact the nodules.\nThe present invention relates to a heat transfer surface having improved cavities to promote nucleate boiling; an improved method for preparing such cavities; and an improved method of nucleate boiling which utilizes the improved cavity structures. The present method for preparing the improved nucleate boiling cavities comprises deforming the metal substrate to form indentations in the metal in the form of non-continuous pits or craters and then electrodepositing a metal on the periphery of the pits or craters at a high current density and then at a normal plating current, from a special bath, to provide honeycomb-like surfaces which optimize nucleate boiling augmentation for particular fluids, especially halocarbon refrigerants utilized in refrigerator systems.\nThe cavity diameters in the plated coatings are directly related to the size of the pits or craters formed on the metal substrate. In the case of sandblasting, the metal deformation is made by grits having a mesh size of 10 to 200. In the case of deforming by knurling, the cavity diameter is determined by the size and number per square inch of the projections on the knurling tool. A pre-plated surface having from 100 to 350 (10,000 to 122,500 per square inch) knurls per inch is preferred in making the present heat exchange surface. It has been found that a knurling tool having 240 knurls per square inch provides pits or craters which, when electroplated, form excellent nucleate boiling cavities in the heated wall surface of a heat exchanger using halocarbon refrigerants such as R-12.\nWhen the heat exchange surface is deformed by sandblasting, silicon carbide of 10-200 mesh, glass beads of 10-200 mesh and ground slag of 60 mesh have been used; silicon carbide of 60 mesh size is preferred for use in R-12."}
{"text": "This invention relates in general to coking devices and in particular to a new and useful apparatus for generating superheated high pressure steam during a dry process of cooling the coke and to an apparatus for carrying out the method.\nA method and apparatus for dry cooling incandescent coke are known from German patent application No. P 33 32 702. In this prior art method, the heat from the hot coke is partly transferred in the cooling chamber to a circulated cooling gas and used in a following waste heat boiler and, through a heat exchanger, for generating superheated high-pressure steam. The usual superheater, evaporator, and feed water preheater are provided in the waste heat boiler and connected in series. Experience has shown that the tubes in the superheater are exposed to extreme conditions. On their outside, they are in contact with the hot circulated gases having a temperature in excess of 1,000.degree. C., and inside, they are in contact with the superheated steam flowing therethrough. Also, the reducing atmosphere of the circulating hot gas require very expensive materials which must withstand these extreme temperature as well as high pressures and very corrosive gases.\nGerman OS No. 30 07 040 discloses a method and apparatus for utilizing the sensible heat of coke in a dry cooling process with a circulated inert gas, where the hot circulated gas is directed from the cooling chamber, upstream of the boiler-heat exchanger system, through a combustion chamber into which a fuel gas and an oxygen containing gas, such as combustion air, are introducible. This additional combustion chamber is intended to ensure a continuous steam generation, if irregularities in the hot coke supply occur. In normal operation of the plant, the entire amount of circulating hot has is directed above into the boiler-exchanger system, and the gas first contacts the coils of the superheater."}
{"text": "1. Field of the Invention\nThis invention relates to a system and method for formatting and displaying numbers, and more particularly, to a system and method for formatting and displaying numbers for use with a hands free telephone in an automobile.\n2. Related Art\nThere are existing systems where a dialed number is parsed and displayed in a designated format. One example of such a system is Lekven et al., U.S. Pat. No. 6,289,226. Another example is Bogart et al., U.S. Pat. No. 5,299,261.\nHowever, these systems are designed to format only telephone numbers and are limited in their ability to format other types of numbers.\nIn addition, neither of these systems or methods provides a system where voice information, as opposed to dialed or typed digits, is received and formatted. With the increasing use of mobile telephones, especially in motor vehicles, hands free operation of a mobile telephone is becoming increasingly important. In some states, it is illegal to operate a mobile telephone unless a hands free system is used.\nThere is currently a need for a system that can receive spoken or dictated information and properly format the information without requiring additional instructions or commands. In other words, there is a need for a system that is truly hands free and can function properly without requiring additional manually typed commands."}
{"text": "Technical Field\nThe present disclosure relates to a vehicle body lower structure of a vehicle such as a car.\nRelated Art\nThe vehicle body lower structure of Patent Literature 1 includes front side members which each have a main body portion on the engine room side and a rear portion extending below a vehicle cabin. The rear portion of the front side member is bent obliquely outward and toward the rear of the vehicle and a rear end of the rear portion is coupled to a rocker rail. A reinforcement member extending obliquely inward and toward the rear of the vehicle is coupled to a vehicle inner side surface of a bent portion of the front side member. The rear portion of the front side member and the reinforcement member extend along a floor panel."}
{"text": "Although originally a speaker unit that forms sound by generating wave in the air, the energy for forming sound cannot be effectively converted to sound in the existing speaker unit. In other words, about half of the energy is converted to vibration of the unit itself without becoming a sound thereby causing distortion of sound and dull vibration, and inhibiting an accurate sound formation, which is the object of the unit.\nIn a speaker system in which the speaker unit is incorporated in an enclosure (speaker box), the enclosure itself similarly vibrates by the vibration of the speaker unit, and such vibration is transmitted to the floor. Generally, the sound becomes difficult to accurately reproduce as the sound becomes lower unless the enclosure is securely fixed to the floor. Thus, if a low sound is emitted at large volume, the entire room may shake by the generated vibration, and the vibration may be transmitted to the adjacent room."}
{"text": "1. Field of the Invention\nThe present invention relates to a multi-layered circuit substrate for a semiconductor device. Further, the present invention relates to a method of manufacturing the multilayered circuit substrate. Furthermore, the present invention relates to a semiconductor device using such a substrate and a method of producing the same. More particularly, the present invention relates to a multilayered circuit substrate for a semiconductor device having a semiconductor element mounting face on which a plurality of semiconductor elements can be mounted being arranged in the plane direction on one side of the multilayered circuit substrate composed by laminating a plurality of conductor patterns. Further, the present invention relates to a method of producing the multilayered circuit substrate, a semiconductor device, and a method of producing the same.\n2. Description of the Related Art\nRecently, there has been provided a semiconductor device referred to as a system-in-package (SIP) in which a plurality of semiconductor elements are mounted on one circuit substrate. This semiconductor device will be referred to as “SIP”, hereinafter.\nA size of the aforementioned SIP is appropriate to be handled as a chip. An example of SIP is shown in FIGS. 7(a) and 7(b). SIP 100 shown in FIG. 7(a) is composed in such a manner that semiconductor elements 104a, 104b, 104c, the functions of which are different from each other, are mounted in the plane direction on a semiconductor element mounting face formed on one face of a piece of multilayered circuit substrate 102.\nAs shown in FIG. 7(b), the multilayered circuit substrate 102 forming the above SIP 100 is composed of resin layers 102a, 102b, 102c, 102d which are made of insulating resin, and the conductor patterns 106, 106, . . . and the via holes 108, 108, . . . are formed being laminated on these resin layers 102a, 102b, 102c, 102d. \nOn one face of the multilayered circuit substrate 102, there are provided connection pads 110, 110, . . . , from which the connecting faces to be connected with the electrode terminals of the semiconductor elements 104a, 104b, 104c are exposed. On the other face of the multilayered circuit substrate 102, there are provided external connection pads 114, 114, . . . , from which the attaching faces on which the solder balls 112, 112, . . . are attached are exposed.\nThe connection pads 110, 110, . . . and the external connection pads 114, 114, . . . are electrically connected with each other by the conductor patterns 106, 106, . . . and the via holes 108, 108, . . . which are formed and laminated on the resin layers 102a, 102b, 102c, 102d. \nOne face and the other face of the above multilayered circuit substrate are covered with the protective films 116, 118 made of solder resist except for the connection pads 110, 110, . . . and the external connection pads 114, 114, . . . .\nIn this connection, the potting resin 120 is filled between the semiconductor element mounting face, which is formed on one face of the multilayered circuit substrate 102, and the semiconductor elements 104a, 104b, 104c mounted on the semiconductor element mounting face.\nWhen SIP 100 shown in FIGS. 7(a) and 7(b) is used, the size of SIP can be reduced to be smaller than the size of a device in which a plurality of semiconductor devices, each semiconductor device having a single semiconductor element, are used. Further, it is possible to reduce a conductor distance between the semiconductor elements 104a, 104b, 104c. Therefore, a signal can be sent and received between the semiconductor elements at high transmission speed.\nHowever, SIP 100 shown in FIGS. 7(a) and 7(b) is composed of a multilayered circuit substrate 102 actually made of resin. Therefore, rigidity of the multilayered circuit substrate 102 is not sufficiently high.\nTherefore, the present inventors made investigations into SIP 200 shown in FIGS. 8(a) and 8(b). As shown in FIG. 8(b), the multilayered circuit substrate 202 composing this SIP 200 is formed by laminating the resin layers 202a, 202b made of insulating resin on which conductor patterns 106, 106 . . . and the via holes 108, 108 . . . are formed. On one face of the multilayered circuit substrate 200, there are provided connection pads 110, 110, . . . , from which the connecting faces to be connected with the electrode terminals of the semiconductor elements 104a, 104b, 104c are exposed.\nOn the other face of the multilayered circuit substrate 202, the metallic plate 204, which is a plate-shaped member having a rigidity higher than that of the multilayered circuit substrate 202, is bonded by the adhesive layer 206 made of insulating resin.\nSince metallic plate 204 is joined to the other face of the multilayered circuit substrate 202, SIP 200 and other electronic parts are electrically connected with each other by a lead frame as shown in FIGS. 8(a) and 8(b). Specifically, an end portion of each inner lead 300, 300, . . . of the lead frame is connected with an exposure face of each external connection pad 208, 208, . . . formed along the outer edge of the multilayered circuit substrate 202 by the wire 302.\nIn this SIP 200, the highly rigid metallic plate 204 is joined to the other face of the multilayered substrate 202. Therefore, rigidity of this SIP 200 is actually enhanced and higher than that of the multilayered circuit substrate 102 made of resin shown in FIGS. 7(a) and 7(b).\nHowever, since the coefficient of thermal expansion of the multilayered circuit substrate 202 actually made of resin and that of the metallic plate 204 are different from each other, stress is generated between them. Cracks tend to be caused on SIP 200 by the thus generated stress.\nSince the multilayered circuit substrate 202 shown in FIGS. 8(a) and 8(b) is successively laminated from the resin layer 202a provided on the metallic plate 204 side. Therefore, on a surface of the resin layer 202b laminated on the resin layer 202a, especially on a surface of the resin layer 202b corresponding to the via hole 108 formed on the resin layer 202a, a concave surface, or indentation, tends to be formed as shown in FIG. 9. Therefore, the surface of the resin layer 202b tends to become irregular.\nWhen the connection pad 110 to be connected with the electrode terminal of the semiconductor element is formed on the irregular face of the resin layer 202b as shown in FIG. 9, an exposed face of the connection pad 110 is formed into an irregular face, following the irregular face of the resin layer 202b. \nIn the case where the exposed face of the connection pad 110 is formed into an irregular face, when a semiconductor element is mounted on a semiconductor element mounting face of the multilayer circuit substrate 202, an electrode terminal of the semiconductor element does not come into contact with the exposed face of the connection pad 110, which causes an imperfect contact and reliability of the finally obtained SIP 200 is deteriorated.\nAlthough the irregular face formed on the resin layer 202b of the multilayer circuit substrate 202 can be flattened by means of polishing, it is necessary to add a polishing process to the conventional manufacturing process of SIP 200, which raises the manufacturing cost of SIP 200. For the above reasons, it is preferable that the semiconductor element mounting face of the multilayered circuit substrate 202 is flattened without adding the polishing process."}
{"text": "Prunus persica. \nIn a continuing effort to improve the quality of shipping fruits, I, the inventor, typically hybridize a large number of peach, nectarine, plum, apricot, and cherry seedlings each year. The present invention relates to a new and distinct variety of nectarine tree, which has been denominated varietally as xe2x80x98SEPTEMBER BRIGHTxe2x80x99. The present variety was hybridized by me in 1992, grown as a seedling on its own root in my greenhouse, and transplanted to a cultivated area of my experimental orchard at Bradford Farms near Le Grand, Calif. in Merced County (San Joaquin Valley). The variety was developed as a first generation cross using xe2x80x98August Redxe2x80x99 (U.S. Plant Pat. No. 6,363) yellow flesh nectarine as the selected seed parent and an unnamed nectarine as the selected pollen parent. Subsequent to origination of the present variety of nectarine tree, I asexually reproduced it by budding and grafting in the experimental orchard described above, and such reproduction of plant and fruit characteristics were true to the original plant in all respects. The reproduction of the variety included the use of xe2x80x98Nemaguardxe2x80x99 rootstock (unpatented) upon which the present variety was compatible and true to type.\nThe present variety is similar to its selected seed parent, xe2x80x98August Redxe2x80x99 (U.S. Plant Pat. No. 6,363) nectarine, by producing nectarines that are globose in shape, firm in texture, mostly red in skin color, and yellow in flesh color, but is distinguished therefrom and an improvement thereon by producing fruit that matures about two weeks later.\nThe present variety is most similar to xe2x80x98September Redxe2x80x99 (U.S. Plant Pat. No. 5,664) nectarine by producing nectarines that are yellow in flesh color and that mature in September, but is distinguished therefrom and an improvement thereon by having a small nonshowy blossom and by producing fruit that is firmer in texture, that has a higher percentage of red skin color, and that is larger in size."}
{"text": "The present invention relates generally to storing information for use in debugging computer software. More particularly, the invention provides a method and apparatus for using a serial bus, such as an IEEE 1394 bus, to remotely create, from a host computer, a snapshot of physical memory of a target computer.\nMany operating systems support xe2x80x9ccrash dumpsxe2x80x9d or xe2x80x9ccore dumps.xe2x80x9d These are typically files that an operating system generates when the operating system recognizes an internal failure. Operating systems typically switch to a minimal mode of operation and attempt to save the operating system\"\"s state information to a file on disk. A software developer or system administrator can subsequently use the saved state information to analyze the operating system failure. Such analysis is often referred to as xe2x80x9cdebugging.xe2x80x9d\nTypically, a crash dump file is generated by the operating system and stored on a disk drive that is local to the computer on which the internal failure occurred. FIG. 2 is a simplified state diagram showing that upon the occurrence of an internal failure, a state transition 202 occurs from normal operation state 200 to minimal mode of operation state 204 in which a crash dump file is stored locally, for instance on a local hard disk.\nUnfortunately, upon the occurrence of an internal failure, for various reasons the operating system may not be able to store the crash dump file locally. For instance, there may be insufficient local storage space available or the target computer\"\"s local hard drive might be off line.\nUnder such circumstance, it is typically desirable to retrieve the crash dump remotely and store it on a remote computer the moment the target computer crashes. Due to the relatively slow rate at which conventional serial buses transfer data, however, transferring the contents of physical memory of a target computer, for instance 64 megabytes of memory, would take a prohibitively long time.\nAccordingly, there is a need for improved techniques for creating a crash dump file much more quickly, while minimizing the amount of operating system state information that is corrupted or otherwise lost, and under circumstances in which a crash dump file cannot be stored on a local drive of the computer that has experienced an internal fault.\nReferring to FIG. 3, techniques for remote debugging of the kernel of an operating system and application programs of a first computer, such as target computer 300, from a second computer, such as host computer 302 are known in the art. In such an environment, a hardware link 304 (e.g., RS-232C), such as a serial connection, is used to transmit debugging commands to and return debugging results from software undergoing test on target computer 300, via serial ports 306 and 308. A small portion of core operating system (xe2x80x9cO.S.xe2x80x9d) 310, typically xe2x80x9cpushesxe2x80x9d data to host debugger 312, thus using processor time on target computer 300 and causing side effects on target computer 300 that would not occur absent debugging-related operation of this small portion of core O.S. code.\nKernel debugging presents technical problems different than those presented by debugging application programs 314, which typically run in user mode. For instance, core operating system 310 typically must be halted for kernel debugging, but core operating system 310 typically is not halted for debugging of user mode application programs.\nFIG. 4 is a flow chart of simplified steps for conventional remote debugging of core operating system 310, which typically runs in kernel mode. At step 400, a user, such as a system administrator, issues a command through host debugger 312 to break execution, in other words stop, suspend, or halt, execution of target computer 300. At step 402, the system administrator may issue well-known debugging commands to step through portions of code, set the contents of target computer memory locations to specific values, and the like. Debugging in this manner is often a very time-consuming process. At step 404, a determination is made regarding whether the system administrator is finished debugging the core operating system 310 of target computer 300. Steps 402 and 404 continue in a loop until the system administrator is finished remotely debugging the kernel of target system 300, at which point host computer debugger 312 releases control of core 0.S. 310, which may resume executing, as indicated at step 406.\nA significant disadvantage of conventional remote kernel debugging in accordance with FIGS. 3 and 4 is that target computer 300 is not operating while the kernel is being remotely debugged. Because debugging is often necessary for certain types of software issues, such as memory leaks, application hangs, and many other non-fatal crashes, and because debugging is typically very time-consuming, there is a need for improved remote kernel debugging techniques that allow debugging to be performed in parallel with operation of a target computer.\nA system and method in accordance with certain inventive principles overcomes the foregoing shortcomings of conventional generation of crash dump files and remote kernel debugging. In accordance with principles of the invention, a two-computer crash dump scenario is provided in which a snapshot of the physical memory of a target computer is extracted and stored by a host computer over a serial bus, which may be IEEE 1394 compliant. In this manner, a crash dump can be transferred quickly and under circumstances in which a conventional single computer crash dump may not be feasible.\nA host computer debugger may remotely issue a command to stop execution of the core operating system of the target computer. Handshake information is then preferably provided from the target computer core operating system to the host computer debugger. The handshake information preferably specifies, among other information, one or more address ranges at which physical memory is present on the target computer, as disclosed in the commonly assigned patent application entitled xe2x80x9cMethod and Apparatus for Remotely Debugging Computer Software Over a Serial Bus,xe2x80x9d application Ser. No. 09/488,015, filed Jan. 20, 2000, which is incorporated herein by reference.\nThe host computer debugger then preferably directly accesses the physical memory of the target computer system over the serial bus in accordance with the handshake information. The host computer then stores the contents of the physical memory of the target computer system.\nExecution of the core operating system of the target system may then be resumed, and the core operating system of the target computer may be debugged in parallel with resumed execution of the target computer core operating system. In this manner, the amount of time that the target system computer is not operational due to debugging is significantly reduced relative to debugging the core operating system of the target computer in real time.\nAdditional features and advantages of the invention will be apparent upon reviewing the following detailed description."}
{"text": "The geometries in integrated circuits have scaled to smaller and smaller dimensions. As the geometries have scaled the alignment tolerance of a pattern level to underlying pattern levels has also become smaller and more critical.\nIn conventional integrated circuit flows, when a pattern such as active, gate, contact, metal, and via is formed on the integrated circuit wafer, alignment marks for alignment in the x-direction and alignment marks for alignment in the y-direction are formed on each of these layers so that subsequent patterns may be aligned using these alignment marks.\nConventional photolithographic alignment is illustrated in FIGS. 1-3. FIG. 1 shows an integrated circuit wafer with multiple integrated circuit chip 20 patterns formed on the wafer. For example, integrated circuit pattern 20 may be the active pattern etched into the integrated circuit wafer. When the active geometries are etched into the integrated circuit wafer active alignment marks 22, 24, 26, and 28 are also etched into the integrated circuit wafer. Typically the alignment marks consist of vertical 40 and horizontal 42 grids as shown in FIG. 3. When another integrated circuit pattern is to be printed on the integrated circuit wafer, it is critical for the new pattern geometries to be properly aligned to the active geometries already etched into the integrated circuit wafer. The reticle 30 illustrated in FIG. 2 may be an integrated circuit gate reticle, an integrated circuit contact reticle, or an integrated circuit implant pattern reticle. A gate reticle is used for illustration. The gate reticle is loaded into a photolithography stepper and the stepper aligns the gate alignment marks 32, 34, 36, and 38 that are on the gate reticle to the active alignment marks 22,24,26,28 that are etched in the integrated circuit wafer. The stepper aligns the gate reticle alignment marks in the x and y directions to the underlying active pattern alignment marks for each of the integrated circuit chips on the wafer. Vertical alignment grids 40 are used for horizontal alignment. Horizontal alignment grids 42 are used for vertical alignment.\nAn alignment tree for a typical integrated circuit aligns all photoresist pattern alignment marks prior to the contact level to the active x and y alignment marks, aligns the contact pattern alignment marks to the gate x and y alignment marks, aligns the first interconnect pattern alignment marks to the contact x and y alignment marks, aligns via pattern alignment marks to the underlying interconnect x and y alignment marks and aligns interconnect pattern alignment marks to the underlying via x and y alignment marks.\nIn addition, in order to pattern geometries with a pitch below about 100 nm while using 193 nm lithography, it has become necessary to use double pattern technology (DPT). Using DPT, one interconnect pattern with a pitch of about 100 nm may be split into two patterns each with a pitch of about 200 nm. The interconnect pattern with a 100 nm pitch may then be formed on an integrated circuit wafer by printing both patterns with the looser pitch into the same layer of photoresist."}
{"text": "It is well known within the electronics industry that lead coplanarity is essential to the successful placement of fine-pitch surface mount devices (SMD's). Non-coplanar leads, i.e. those bent up or down with respect to the body of the component, typically result in opens or poor solder joints on the printed circuit board (PCB). It is desirable to reject any components that do not meet coplanarity requirements prior to placement on the printed circuit board. The present industry standards require that all lead tips be coplanar within a 0.004 inch band, forcing inspection equipment to measure coplanarity with .+-.0.002 inch accuracy. Although some non-contact 3-D inspection techniques are capable of providing this degree of accuracy, a number of system constraints prohibit their implementation for this application. These constraints include but are not limited to:\n(1) large part sizes,\n(2) establishment of a reference (datum) plane,\n(3) mechanical space and packaging limitations within the pick and place machine, and\n(4) real-time inspection to maintain high pick and place throughput rates.\nOther techniques, particularly \"implied\" techniques based around passive triangulation (i.e. estimation of lead height via measured lead length), have severely limited depth resolution. They may also have other drawbacks such as confusion between lead length error and vertical lead displacement as well as ambiguous polarity information (the inability to differentiate between leads bent up or down with respect to the plane of the part).\nIt is accordingly an object of the present invention to overcome these constraints while providing the necessary degree of accuracy for proper inspection of lead coplanarity. The optical setup inherent in the disclosed system provides a true, direct measurement of lead height. It produces sufficient depth resolution to measure lead coplanarity to an accuracy of .+-.0.002 inch using conventional sensors and image processing hardware.\nAn advantage of the disclosed system is its ability to handle a range of SMD components including PQFP, SOIC, PLCC, SOJ, and TapePak with pitch values down to 15 mils. It will accommodate 25-mil pitch parts up to 1.25 inches square in a single field of view. Parts with a larger number of leads or finer leads can also be inspected by breaking up the image into multiple fields of view, typically two or four. If necessary, large fine-pitch devices may be inspected through the use of high resolution image sensors such as line scan or non-RS-170 area cameras, though at present this requires additional complexity for image acquisition.\nOther objects and advantages of the present invention will become apparent from the following portion of the specification and from the following drawings which illustrate in accordance with the mandate of the patent statutes a presently preferred embodiment incorporating the principles of the invention."}
{"text": "My new variety of lily plant originated as a seedling which first flowered in 1988 in Boring, Ore. The breeding efforts had as their objective the production of colored, fragrant hybrids between Lilium longiflorum (the Easter lily) and Asiatic lilies, suited to forcing into flower out of season, heretofore unknown in the lily breeding art.\nI achieved the desired objective by intercrossing a slightly fertile embryo-cultured seedling from (Lithium longiflorum `Ace`.times.`Sterling Star`) with `Connecticut King.` The seed parent was the diploid form of the Dutch cultivar `Longistar` (unpatented), which had limited commercial distribution in the Netherlands and the United States. This cultivar does not normally produce fertile seed; it can be used in breeding only if its embryos are artificially cultured. The pollen parent, `Connecticut King,` although unpatented, has been extensively cultivated as a commercial cut-flower variety in the United States and the Netherlands. The embryo that produced `Yellow Fashion` was artificially cultured as well.\nThe flowers of my new lily are characterized by an very large, fragrant, yellow upfacing to slightly outfacing flowers with an \"open bell\" form with slightly clawed tepals, quite narrow at the base and broadest just below the midpoint of the bud and flower, and with tepal tips flattening but recurving only slightly. This unique flower form is intermediate between the trumpet-form of Lilium longiflorum and the \"flat-faced\" Asiatic hybrid form with conspicuously recurved tepal tips. The new clone possesses unusually strong, stout stems. In addition, the clone possesses to a high degree desirable characteristics of hybrid vigor. The clone is a good grower and propagator, as observed at Boring, Ore.\nMy new variety of lily plant has been asexually reproduced by me and under my direction at Boring, Ore. Successive generations produced by natural propagation from bulblets, by bulb scale propagation, and by tissue culturing from bulb scale explants have demonstrated that the novel and distinctive characteristics of my new variety are fixed and hold true under asexual propagation from generation to generation."}
{"text": "The present invention relates to a valve actuating apparatus for an internal combustion engine. More particularly, the invention relates to a valve actuating apparatus which is capable of selectively disabling valves of the engine. This invention finds particular application in motorcycle engines.\nTo obtain an improved operating efficiency over a wide range of engine speeds, engines have been developed having multiple intake and exhaust valves for each cylinder. For instance, in an engine having two intake and two exhaust valves for each cylinder, for low and medium engine speeds, one of the intake and one of the exhaust valves may be disabled, while all four valves are enabled for higher engine speeds. Of course, a valve actuating apparatus must be provided to selectively actuate the valves.\nFIG. 1 shows a cross section of a part of a conventional valve actuating apparatus utilizing an oil pressure actuator 5 for selectively disabling the intake and exhaust valves. In this valve actuating apparatus, a pair of rocker arms 3 and 4 are adjacently mounted on a rocker arm shaft 2, the latter being fixedly mounted to a cylinder head 1. One of the rocker arms 3 is directly driven by a cam shaft, and the other rocker arm 4 is selectively engageable with the rocker arm 3 to thereby selectively enable the valve associated with the rocker arm 4. The actuator 5 is a reciprocation type, constructed such that a shaft pin 5a slidably mounted in the rocker arm 3 is pushed by pressurized oil to protrude into the rocker arm 4. A pressure oil chamber 5b receiving the shift pin 5a is communicated through an oil path 5c with an oil path 5d provided in the rocker arm shaft 2. An engaging hole 5e sized receiving the shift pin 5a therein is provided at a position of the rocker arm 4 opposing the pressure oil chamber 5b. There is provided in the engaging hole 5e a return pin 5f for urging the shift pin 5a back into the pressure oil chamber 5b and an elastic member 5g for urging the return pin 5f toward the rocker arm 3.\nIn operation, when the engine is driven at a high speed, operating oil at a predetermined pressure is supplied into the oil path 5d from an oil supply apparatus (not shown) to thereby energize the actuator and couple the rocker arms 3 and 4.\nIn most conventional actuator systems, oil paths are provided in members adjacent to the actuator (the rocker arm 3 and the rocker arm shaft 2, etc. in the abovedescribed actuator) to provide paths for supplying operating oil as lubricating oil for other engine components. Specifically, interlinked portions of the members forming the oil paths and the sliding portion of the actuator are constructed so as to allow the operating oil to be leaked slightly therefrom.\nIf the engine has been operated at a low speed for a long period, and hence the actuator is not supplied with operating oil for a long period, air tends to leak into the oil paths from the sliding and the interlinked portions. If the air leaked into the oil paths remains until the actuator is activated, the responsiveness of the actuator may be low. Namely, in the above-described valve actuating apparatus, for example, the moving speed of the shift pin 5a is made lower so that the edge portion of the shift pin 5a is likely to collide with the edge portion of the inlet of the engaging hole 5e thus making it difficult to smoothly activate the actuator to couple the rocker arm 4 to the rocker arm 3.\nTo overcome this drawback, it has been proposed to provide an oil pump for exclusive use by the actuator adjacent to the actuator so as to shorten the time required for the operating oil to reach the actuator from the oil pump. However, in this case, the size of the engine is unavoidably increased due to the provision of the oil pump near the actuator, and further a transmission mechanism for exclusive use for driving the oil pump is required, thereby making the oil supply apparatus expensive.\nA typical example of the conventional valve actuating apparatus installed on an engine is shown in FIG. 3. There is provided, in a cylinder head a at the top of the engine disposed almost at the center of a body having a front and a rear wheel, a valve actuator d operated by pressurized oil and having an oil pressure chamber c. Operating oil is supplied to the oil pressure chamber c from a control valve b, whereby the intake and exhaust valves e of the engine are selectively activated. The conventional control valve includes a spool valve j for selectively opening a control path f, an oil supply path g and an oil return path h provided in a valve housing i having a control path communicating with the oil pressure chamber c. The oil supply path g communicates with an oil pressure source and the oil exhaust path h communicates with an oil ejecting opening.\nIn this type of control valve, the height of an opening end of the oil exhaust path h, that is, the height of the oil ejecting opening, is of particular importance. If the oil ejecting opening is disposed relatively lower, the oil in the path h is likely to undesirably flow out through the oil ejecting opening, while if it is disposed relatively higher, the resistance of the oil path from the control path to the oil ejecting opening becomes higher so that the responsiveness of the control apparatus is lowered. Further, the control valve must be positioned by taking into consideration the tilt angle thereof when the motorcycle is held up by its kickstand.\nThe control valve b is generally disposed along a horizontal line parallel to the crankshaft of the engine in the valve casing i as shown in FIG. 3. Thus, the control valve b is subjected to the vibration of the engine, which may cause the state of the control valve b to change states unintentionally. Namely, for example, a four-cycle, four-cylinder in-line engine generally vibrates not only in a vertical direction, but also in a horizontal direction (the direction parallel to the crankshaft) with an amplitude almost half that in the vertical direction. In order to prevent such erroneous changeover operations of the intake and exhaust valves, it has been unfavorably required to increase the elasticity of a spring or the magnetic force of an operating solenoid associated with the control valve."}
{"text": "With development of display technology, liquid crystal displays have become most commonly used display devices.\nIn the process of manufacturing a liquid crystal display, the cell process, as an important procedure, comprises not only aligning and then assembling an array substrate and a color filter substrate, but also irradiating the frame glue filled between the array substrate and the color filter substrate by using ultraviolet light so as to solidify the glue. Since the electrical performance of the thin film transistor (TFT) provided on the array substrate will be affected after ultraviolet irradiation, those parts where the TFT and the like are located should be covered with an ultraviolet mask plate during the irradiation, but leaving the frame glue exposed.\nAn existing ultraviolet mask plate, comprising a glass substrate and a layer of shielding metal formed on the glass substrate, is also prepared by using a mask plate through a patterning process. Since the ultraviolet mask plate can be reused, the usage rate of the mask plate for manufacturing the ultraviolet mask plate is rather low. Hence, the cost of manufacturing the ultraviolet mask plate is relatively high in the prior art."}
{"text": "A router is a device that forwards traffic between networks. Routers use headers and a forwarding table to determine where packets go, and they use messaging such as the Border Gateway Protocol (BGP) to communicate with each other and configure the best route between any two hosts.\nConventional routers include a packet forwarding component and a route processor (RP). The RP determines and controls configuration, security, accounting, debugging, and network management processes of the packet-forwarding component. Examples of route processors include the Cisco® CSC/3 and the CSC/4 route processors. Route processors are typically implemented on an electronic printed circuit board.\nThe RP typically has relatively more processing and resource requirements than the packet-forwarding component. As the traffic between networks increases, router capacity must also increase. This places an even greater processing demand on the RP. The RP can, therefore, become a “bottleneck” as processing in the router is slowed by inadequate processing capability within the RP.\nOne solution for meeting the increased processing demands on the RP is to distribute the RP function across two or more loosely coupled processors. A problem with this approach is how to direct locally addressed packets to the appropriate RP out of the two or more distributed route processors. Another problem is how to reassemble fragments in an RP distributed across two or more processors.\nFor the reasons stated above, and for other reasons stated below which will become apparent to those skilled in the art upon reading and understanding the present specification, there is a need in the art for a system and method for routing packets within routers having distributed route processors."}
{"text": "Wearable personal ornamental items are widely popular in cultures throughout the world. Such items include, for instance but not limited to, watches, bracelets, and rings. Rings are widely worn ornaments used to decorate, for example, fingers and toes. Rings may be designed to fit a finger differently. For instance, “regular fit” bands include a flat interior. U.S. Pat. No. 6,166,053, to Seibenberg, reports “comfort fit” bands including a curved interior as reported in and shown in FIG. 2 therein. A ring interior including an annular circumferential groove is reported by U.S. Pat. No. 6,701,618, to Gefen.\nMethods of fabricating rings and other jewelry items are known to those skilled in the art, and are set forth, for instance, in “The AJM Guide to Lost-Wax Casting” (MJSA/AJM Press); Mulock, J., “Cast-A-Way: The Complete Lost Wax Process for the Amateur or Professional Jeweler”; Romanoff, P., “The Complete Handbook of Centrifugal Casting”; and Romanoff, P., “The Art and Science of Centrifugal Casting Metal.”\nUse of gemstones, enameling, engraving, and finishing techniques in jewelry has also been described. Gemstone-setting techniques are described in Wooding, R. R., “Diamond Setting Manual: Procedures and Techniques”; Wooding, R. R., “Diamond-Setting: The Professional Approach”; Wooding, R. R., “Bead Setting Diamonds with Pave Applications”; and Wooding, R. R., “Channel Setting Diamonds with Illustrated Procedures.” Jewelry finishing is described, for example, in Hoch, “Tumble Finishing for Handmade Jewelry.” Engraving of jewelry is described, for example, in Meek, J. B., “The Art of Engraving.” Electroplating and electroforming techniques are described, for example, in “Midas® Plating Guide,” Third Edition, produced by Rio Grande. Enameling techniques are described, for example, in McGrath, J., “First Steps in Enameling”; Werge-Hartley, J., “Enamelling on Precious Metals”; Cohen, K. L., “The Art of Fine Enameling”; and Matthews, G. L., “Enamels, Enameling, Enamelists.” Jewelry may also be constructed by die-striking or stamping, as discussed in McGrath, J., “The Encyclopedia of Jewelry-Making Techniques: A Comprehensive Visual Guide to Traditional and Contemporary Techniques.”\nA secure fit is often desirable in a ring to prevent inadvertent loss of the ring. Preferably the fit of the ring is not so secure that the wearer experiences discomfort. Other personal ornament items may be designed to effect different fit characteristics including a secure fit similar to that desirable in a ring. For example, bracelets or watch bands may be constructed to have a snug, or tight, fit.\nAlthough a tight fit may be desirable in a jewelry item, such a fit may cause issues related to ease of use of the jewelry item. For instance, a tightly-fitting ring may be difficult to put on or to remove. Once placed, for example, on a digit or around a wrist or ankle, a tightly-fitting jewelry item may also not be comfortable to the wearer. Previous attempted solutions for increasing the comfort of a tightly-fitting jewelry item, such as those reported by patents mentioned above, may not provide a level of comfort that is entirely satisfactory.\nAll patents and documents referred to herein are hereby incorporated by reference as if rewritten herein in their entirety."}
{"text": "1. Field of the Invention\nThe present invention relates to the production of poly(ethylene-2,6-naphthalene dicarboxylates).\n2. Discussion of the Background\nPoly(ethylene-2,6-naphthalene dicarboxylates) are produced by the polyesterification of dimethyl-2,6-naphthalene dicarboxylate and ethylene glycol in the presence of polycondensation and polyesterification catalyst systems.\nIn the preparation of dimethyl-2,6-naphthalene dicarboxylate, significant amounts of methyl hydrogen 2,6-naphthalene dicarboxylate are also produced. Consequently, during a typical preparation of poly(ethylene-2,6-naphthalene dicarboxylate) from dimethyl-2,6-naphthalene dicarboxylate, the presence of methyl hydrogen 2,6-naphthalene dicarboxylate inhibits the polycondensation reaction. This inhibition produces a final product that has a lower than desirable inherent viscosity.\nThere is thus a need for producing a high quality poly(ethylene-2,6-naphthalene dicarboxylate) product from a dimethyl-2,6-naphthalene dicarboxylate starting material contaminated with methyl hydrogen 2,6-naphthalene dicarboxylate."}
{"text": "New communication services and the use of complex waveforms have created a demand for highly linearized power amplifiers. Deviations from linearity are manifest by spectral distortion components in the output of these amplifiers—that is, undesired energy, not contained in the original signal, inside or outside the frequency band of interest. Linearization techniques seek to reduce these distortion components, restoring signal integrity and allowing an amplifier to operate at its best spectral and power efficiency for the specific application.\nTo date, a number of linearization approaches have been attempted. Most of these techniques, however, are useful only in circumstances where the input signals occupy low-percentage bandwidths (typically 10% or less) or do not change frequency rapidly over time.\nFor example, existing digital pre-distortion (DPD) techniques, which use non-linear transfer characteristics of an amplifier in order to generate inverse distortion products that can be supplied along with input signals to the amplifier, require time to converge and must have access to the baseband or digital intermediate frequency (IF) signal components (I and Q or IFI and IFQ). Moreover, digital signal processors, ASICS or FPGAs used in DPD-based designs tend to be limited in bandwidth (due to their clock frequencies) and processing capacity.\nCurrent analog pre-distortion techniques are likewise limited in bandwidth and require adjustment when the input signal center frequency is changed. Adaptive linearization can eliminate the limitations to the change in the center frequency of these designs, but tend to be slow to respond to changes in system characteristics and can be very complex to implement.\nFeed-forward linearization techniques are also complex to implement (often making them unsuitable for use with existing amplifiers), are limited in bandwidth and often do not provide significant gains in efficiency because of the relatively high power consumption of the error or auxiliary amplifier. Other conventional linearization solutions, such as power back-off schemes, envelope elimination and restoration designs, linear amplification using nonlinear components, and Cartesian feedback designs all have their own limitations as well. In addition, the increasing complexity of semiconductor-based power amplifiers has memory effect or inertia and therefore limits the use of converging linearization techniques in frequency hopping applications due to the reaction time to correct."}
{"text": "In the past, removal of the access panel or cover to the chassis of an electronic device such as a computer, or the like, was cumbersome or time consuming. In most applications, the access panel is mounted to the chassis using fasteners such as thumb screws, or the like, which must be unscrewed each time the access panel is removed from the chassis. The access panel may then be lifted or slid from the chassis to provide access to components located therein. More recently, quick release latching systems have been provided for securing the access panel to the chassis of the computer. Typically, such latching systems also push the access panel open when the latching mechanism is defeated, making the access panel easier to remove by users.\nIn modern computer manufacturing facilities, vacuum lifting is often used for lifting the computer to place the computer in packaging prior to delivery. However, when vacuum lifting is employed with a computer chassis that utilizes a quick release latching system that pushes the access panel of the computer open, the access panel can be inadvertently pulled from the chassis of the computer, causing the computer to be dropped and possibly damaged. Generally, this problem occurs because the devices that defeat the latching mechanism are actuated by forces that are applied normal to the access panel during vacuum lifting. The latching mechanism then works in concert with the forces applied during the vacuum lifting process to separate the access panel from the chassis.\nConsequently, it would be desirable to provide a quick release latching mechanism that allows the access panel of a computer or similar electronic device to be easily removed from its chassis, but which is not susceptible to inadvertent actuation during processes, such as vacuum lifting, that apply forces to the access panel."}
{"text": "1. Field of the Invention\nThe present invention relates to an image signal compressing device adapted for use in a facsimile machine and an electronic filing machine.\n2. Description of the Prior Art\nRecently, a device has been developed which, when used in a facsimile machine or an electronic filing machine, operates to convert an image of intermediate gradation into a binary image signal according to a dither process.\nIt has, however, been found that, when the binary image signal, converted by the above described device, is compressed according to an MR code or a modified Hoffman (MH), which is a coding (compressing) system standardized in the facsimile machines, the compression ratio tends to be considerably lowered."}
{"text": "Hydroxamic acids are an important class of organic molecules playing a key role in many biologically relevant interactions. Inhibition of matrix metalloproteinases (MMPs) (see Greenwald, R. A.; Golub, L. M., Eds.; Inhibitioin of Matrix Metalloproteinases: Therapeutical Potential; New York Academy of Sciences, vol. 732 (1994), and Rockwell, A.; Melden, M.; Copeland, R. A.; Hardman, K.; Decicco, C. P.; DeGrado, W. F.; J. Am. Chem. Soc., vol. 118 (1996), p. 10337 and references therein) or a unique deacetylase of lipid A biosynthesis (see Onishi, R. H.; Pelak, B. A.; Gerckens, L. S.; Silver, L. L.; Kahan, F. M.; Chen, M-H.; Patchett, A. A.; Galloway, S. M.; Hyland, S. A.; Anderson, M. S.; Raetz C. R. H.; Science, vol. 274 (1996), p. 980) testify to the significance of this class of compounds.\nSynthesis on solid support is a crucial technology for combinatorial chemistry efforts. It allows for easy automation of processes and convenient handling of polar molecules throughout the synthetic protocol. It also provides a reliable method for the preparation of mixtures of compounds (split-mix synthesis). The solid supported synthesis of hydroxamic acids, based on Wang, Sasrin, or the 2-chlorotrityl O-hydroxylamine bound resin have recently been reported. (See Richter, L. S.; Desai, M. C.; Tetrahedron Lett., vol. 38 (1997), p. 321; Gordeev, M. F.; Hui, H. C.; Gordon, E. M.; Patel, D. V.; Tetrahedron Lett., vol. 38 (1997), p. 1729; Mellor, S. L.; McGuire, C.; Chan, W. C.; Tetrahedron Lett., vol. 38 (1997), p. 3311; and Bauer, U.; Ho, W-B.; Koskinen, A. M. P.; Tetrahedron Lett., vol. 38 (1997), p. 7233.) Another alternative synthetic approach, based on resin N-linked hydroxylamine has been reported. (See Ngu, K.; Patel, D. V.; J. Org. Chem., vol. 62 (1997), p. 7088.)\nWe have found two major limitations of the reported methods for making hydroxamic acids using solid-support resins. First, during side chain modifications we sometimes observed undesired functionalization of the nitrogen which ultimately is part of the hydroxylamino moiety of the hydroxamic acid: Resin-ONHCOR. Second, none of the above reported methods allows application of acid labile protecting groups (e.g. Boc) during side-chain synthesis."}
{"text": "Subsea is a well-known term used to refer to equipment, methods and technology used at underwater locations, and particularly at locations on or close to the seabed for various applications, and in particular to explore, drill and develop oil and gas fields that exist below the ocean floors. The ocean floor or seabed at which subsea apparatus is used may be at a relatively shallow location (shallow in this context being at depths less than 1,000 feet) or at a deep water location (generally meaning water depths in the range 1,000 to 10,000 feet, or even deeper).\nIt will be appreciated that a wide variety of types of subsea apparatus (equipment) is used in such applications. This variety includes subsea apparatus providing a control function (for example comprising controllable valves to regulate the flow of fluids to or from a well) and subsea apparatus providing a connection function (such as providing an electrical or optical control or monitoring link between two pieces of apparatus, or providing a fluid conduit from one subsea location to another, for example to convey fluids to a well for injection into that well, or to convey production fluids from a well). A term commonly used to refer to one type of such subsea interconnecting apparatus is “umbilical”, and a subsea umbilical will typically comprise a flexible intermediate portion containing at least one electrical cable, optical wave guide, or fluid conduit, with connection means provided generally at each end for connecting the umbilical to other apparatus. For example, umbilicals are known which comprise a plurality of electrical conductors (e.g. wires or cables) with a corresponding plurality of electrical contacts provided in connectors at each end of the umbilical. There are circumstances in which such an umbilical, or indeed another piece of subsea apparatus, is conveyed to a subsea (underwater) location (this operation is also referred to in the art as “deploying” the apparatus) and is required to remain at that location for some time before it is installed (i.e. connected to some other piece of subsea equipment). For example, an umbilical for providing fluid, electrical and perhaps optical connection to a well head may first be laid on the seabed (deployed). Then, an end of the umbilical may be placed into a mechanical receptacle on the well head to provide the fluid connection (and in this condition the umbilical is partially installed). In certain known arrangements the electrical and/or optical connection between the umbilical and well head is provided by connecting another piece of subsea apparatus, known as a “jumper”, between a connector on the umbilical and a corresponding connector on the well head. Thus the jumper is another form of subsea interconnect apparatus or device, typical shorter than apparatus referred to as umbilicals, but similarly providing an interconnecting function. When the jumper is connected, the umbilical is then fully installed to the well head. Clearly, although a subsea umbilical or other subsea apparatus can be tested before it is deployed (i.e. before it is conveyed to its subsea location), degradation or damage to it may occur while it is being transported to its deployment location and indeed while it is waiting at that location before installation occurs. For example, if the connectors of subsea apparatus are not protected in some way at the subsea location before installation, then they can in general become contaminated with silt or other material, and any electrical contacts they contain could suffer corrosion. Such degradation is problematic. One attempt to solve these problems has been to fit a so-called “dummy” connector to the connector of subsea apparatus such as an umbilical, the dummy connector mating with the apparatus's connector in the same way as would a corresponding connector of apparatus to which the umbilical is to be connected. Certain known arrangements are such that when the connector and dummy connector are mated at least one seal is made which prevents or at least inhibits the ingress of contaminants into the mated arrangement and prevents seawater from making contact with any electrical contact of the umbilical's connector. Certain dummy connectors are adapted to provide protection to other parts of the connection interface (i.e. not just to electrical contacts). The dummy connector typically remains mated with the connector of the subsea apparatus until shortly before that apparatus is to be installed. The dummy connector is then detached and the subsea apparatus is connected in the required configuration. Although this technique offers some advantages, there is still the possibility that the subsea apparatus could have been damaged or degraded in some way while the dummy connector was in place. For example, although the dummy connector may have protected the connector of an umbilical, the umbilical may have suffered damage elsewhere along its length whilst in situ or waiting for deployment. Only after removal of the dummy connector, and connection of the umbilical to other subsea apparatus may this damage or degradation become apparent, when a problem is detected with the operation of the subsea system as a whole.\nIt is therefore an object of certain embodiments of the invention to provide subsea assemblies, subsea test apparatus and methods of handling subsea apparatus which overcome, at least partially, one or more of the problems associated with the prior art."}
{"text": "As is generally known in the field of semiconductor memory devices and other semiconductor integrated circuits, it is often required to generate internally voltages that are greater than an external voltages, also known as off-chip power supply voltages. For example, it is known in flash electrically erasable, programmable read-only memories (EEPROMs) that a first high voltage of about +5V is needed to be produced for the read mode of operation of memory cells. Also, a second high voltage of about +10V is needed to be produced for the program mode of operation of the flash memory cells. To meet this requirement, the semiconductor memories also generally include one or more internal voltage boosting circuits for generating output signals boosted to be higher than an external supply voltage.\nIn FIG. 1, there is shown a simplified block diagram of a conventional voltage boosting circuit 10 for generating a word line voltage VPP, which is passed to appropriate word lines WL0 through WLn in a memory cell array via respective row decoder circuits 20. The voltage boosting circuit 10 includes a high voltage charge pump circuit 14 which is used during the program mode of operation of the semiconductor memory devices. The charge pump circuit 14 serves to charge up the word line voltage VPP at the internal word line supply node N1 to approximately 10V when the enable signal ENVPP is high. The boosting circuit 10 also includes a precharge logic circuit 12 which functions to precharge the word line voltage VPP to the external power supply voltage VCC, before the node N1 is boosted or pumped to a high voltage. The precharge logic circuit 12 will also maintain the supply voltage VPP equal to the external power supply voltage VCC in other modes of operation via the PMOS transistor P1. The capacitor C.sub.L represents the capacitive loading of the row decoder circuits 20, plus all of the parasitic capacitance associated with the line 2 connected to the word line supply node N1.\nWith reference to FIGS. 2(a) through 2(e), the operation of the voltage boosting circuit 10 depicted in FIG. 1 will now be explained. Initially, it is assumed that prior to time t1 the enable signal ENVPP and the precharge bar signal PRECHARGEB are both low (VSS), and that kick bar signal KICKB is high (VCC). Thus, the output of the NAND gate 18 on the line 3, defining the kick voltage VKICK, will be low (VSS). Further, the word line voltage VPP at the supply node N1 will be at the power supply voltage VCC, since the PMOS transistor P1 will be turned on.\nWhen the kick bar signal KICKB makes a high-to-low transition at the time t1 as shown in FIG. 2(a), the precharge bar signal PRECHARGEB in FIG. 2 will be pulled to the power supply voltage VCC at time t2. At time t3, the kick voltage VKICK will begin making a low-to-high transition, as depicted in FIG. 2(b). This will, in turn, cause the word line voltage VPP at the supply node N1 to be raised or boosted by the kick signal VKICK via the booster capacitor C.sub.BOOST, from an initial value of VPPinit=VCC, to a boosted level VPP.sub.BOOST, starting at the time t3 in FIG. 2(d). As a result, the precharge bar signal PRECHARGEB will be further raised to the boosted level at time t4 as shown in FIG. 2(c).\nThis boosted level of the word line voltage VPP.sub.BOOST can be calculated from the following equation: ##EQU1## Where: EQU VKICK=Power supply potential VCC (2) EQU VPP.sub.init =Power supply potential VCC (3)\nand C.sub.BOOST, C.sub.L are respectively the capacitances of capacitors C.sub.BOOST, C.sub.L. By substituting the equations (2) and (3) into the equation (1) and simplifying, there is given: ##EQU2##\nFrom above equation (4), it can be seen that if the required maximum for the word line voltage VPP is less than the level that can be generated by the voltage boosting circuit 10 of FIG. 1, then the conventional boosting circuit 10 will be able to operate adequately. Unfortunately, as previously pointed out, the word line voltage VPP is typically required to be pumped up to approximately +10V for the program mode of operation of the flash memory cells. Therefore, under these circumstances, an additional means is required for producing this higher voltage, such as a second higher voltage charge pump.\nHowever, the use of the conventional voltage boosting circuit 10 for this purpose is impractical, since the total capacitance that is seen by the high voltage charge pump 14 is very high, which is because the boost capacitor 15 is constantly connected to the kick signal VKICK, which is at a value of either VCC or VSS. As a result, this conventional voltage boosting circuit 10 suffers from the disadvantage that either the charge pump 14 would require a much longer time in order to raise the voltage VPP to the necessary programming level, or the size of the charge pump 14 must be increased, so as to raise the voltage VPP to the necessary programming level in the same amount of time. The first way would increase the programming time, and the second way would increase the cost and complexity of the overall memory device. Thus, neither of these choices offers a satisfactory solution.\nMore specifically, this problem of the conventional circuit 10 having a high total capacitance can be illustrated by referring back to equation (2) above. If we let VPPmin be equal to the minimum VPP supply voltage required for the read mode of operation for a flash memory cell, and let VCCmin be equal to the minimum VCC power supply voltage of the flash memory cell, then we have through substitution: ##EQU3##\nFurther, by assuming VPP.sub.min to be equal to +4.2V and VCC.sub.min to be equal to +2.5V, there is given: ##EQU4##\nBy solving for the booster capacitor C.sub.BOOST, there is given: C.sub.BOOST =2.125C.sub.L. Therefore, the total capacitance C.sub.TOTAL has been shown to have increased due to the booster capacitor: C.sub.TOTAL =C.sub.BOOST +C.sub.L or =3.125C.sub.L.\nAccordingly, it would be desirable to provide an improved voltage booster circuit which overcomes the disadvantages of the conventional case, so as to be capable of operating efficiently and effectively in a low supply voltage environment. It would be expedient that the voltage booster circuit be operated effectively so as to produce a significant reduction in power dissipation."}
{"text": "The present invention is directed toward optical communication system. In particular, optical communication systems providing dispersion slope compensation.\nOptical signals transmitted in a fiber optic communication system typically constitute a series of pulses of digital information. Although the pulses are usually at a single nominal wavelength, each pulse is actually composed of different spectral components. These spectral components propagate through the transmission fiber at different speeds with higher frequency components traveling slower than lower frequency components. This effect, known as xe2x80x9cchromatic dispersionxe2x80x9d, can result in spectral components of one pulse arriving at a receiver at substantially the same time as a succeeding pulse, thereby causing degraded receiver sensitivity. Chromatic dispersion becomes increasingly pronounced at higher bit rates, e.g. 10 Gbit/sec. such as those associated with synchronous optical network (SONET) OC192 transmission speeds. Unique fibers have been developed to offset or compensate chromatic dispersion. These fibers, referred to as dispersion compensated fiber or DCF, are commercially available from Corning Inc. and Lucent Technologies, for example. DCF fibers are typically specified with a dispersion coefficient having an opposite sign to that of the transmission fiber. The net dispersion experienced by the transmitted optical signal is thus the sum of the dispersion accumulated through the transmission fiber and the DCF. If the product of dispersion coefficient times length in the DCF is equal in magnitude and opposite in sign to that of the transmission fiber, the net dispersion is zero, and the dispersion is said to be compensated.\nDCF typically provides compensation at a particular wavelength. As a result, single wavelength optical transmission systems incorporating DCF and operating at the specified wavelength can have relatively low error rates, even at OC192 transmission speeds.\nIn order to further increase fiber capacity, wavelength division multiplexed (WDM) systems have been developed for carrying multiple wavelengths on a single fiber. At OC192 rates, each wavelength or channel of a WDM system must be dispersion compensated. Typically, the dispersion coefficient of the transmission fiber, as well as that associated with DCF, varies with wavelength; the relationship of dispersion and wavelength being referred to xe2x80x9cdispersion slopexe2x80x9d. Since the dispersion slope associated with the transmission fiber and DCF may not be identical, not all channels may be adequately dispersion compensated.\nTheoretically, one approach to providing spectrally uniform dispersion compensation involves fabricating DCF uniquely tailored to have a slope of the same magnitude, but opposite sign, as the transmission fiber. As a practical matter, however, such DCF would be difficult to fabricate with precisely the required slope and magnitude.\nAlternatively, since the amount of dispersion compensation depends in part on the length of the DCF, varying lengths of DCF can be provided in a WDM system downstream from the point where individual wavelengths have been separated. Each channel can then be compensated individually prior to detection at a photodetector. An additional segment of DCF, however, must be provided for each channel, thereby complicating system design and increasing costs.\nConsistent with the present invention, an optical device is provided comprising an optical splitter having an input port and a plurality of output ports, said input port being configured to receive a plurality of optical signals, each at a respective one of a plurality of wavelengths. The optical signals are output from each of the plurality of output ports, but in attenuated form. The optical device further includes a plurality of in-fiber Bragg gratings, each of which being coupled to a respective one of the plurality of output ports. The in-fiber Bragg gratings are configured to reflect a selected one of the plurality of optical signals and compensate for a dispersion associated with the selected one of the plurality of optical signals."}
{"text": "There are many reasons to encode an inaudible message in audio data and many groups would like to have access to such technology. A group with such an interest is the group of copyright owners. Copyright owners would like such an encoding technique to facilitate copyright enforcement and protection. Copyright enforcement would be facilitated by encoding pieces of copyrighted works with a watermark to provide ownership information for copyright enforcement. Alternatively, the copyrights of a work may be protected by a copy protection scheme, e.g. encryption keys encoded onto the audio data, which would prevent unauthorized use of the protected matter.\nAnother group with an interest in using inaudible messages encoded into audio data would be the group of audio listeners. The encoding would provide listeners with useful information about the programs they are listening to without affecting the audio experience. For example, the names of the performers, the name of the performance, or the name of the broadcaster may be given and relayed to the listener via the listener's receiver.\nStill another group with an interest in the encoding of inaudible messages into audio data would be market researchers who make use of audience estimating techniques, as well as customer loyalty programs, commercial verification functionality and program identification. Inaudible messages encoded into broadcast or recorded audio are particularly useful in implementing such techniques and activities.\nYet still another group with an interest in the encoding of inaudible messages into audio data would be those seeking additional bandwidth to communicate data that is totally unrelated to the audio data. For example, telecommunications companies could utilize the bandwidth to carry their data and/or news organizations could relay real time news such as breaking headlines or stock quotes.\nThere are many other good reasons that other interested groups have for the encoding of inaudible messages into audio data. One problem encounterd in attempting to encode multiple messages inaudibly within audio data is that there is only a limited amount of bandwidth available for this purpose.\nThe limited bandwidth is due to the fact that audio data can only receive a finite amount of energy in the encoding process before the encoding becomes audible. This level of acceptable ancillary data energy in audio data is application dependent. For example, in high fidelity applications such as music distribution or broadcasting, the messages must be keep inaudible. However, in certain other applications such as voice data communication, e.g. cell phone communications, the constraints on the amount of acceptable ancillary data energy in the audio data are less rigorous. The bandwidth limitations due to these constraints are further restricted by the administrative load imposed by error detection and correction data, marker data, sync data, address data and the like.\nA further problem arises in applications requiring the encoding of one or more messages in audio data that is already encoded with another message. This is desired in certain broadcast and recording applications, such as audience measurement, commercial and network clearance, and content identification. It has been proposed to reserve different respective time intervals along the time base of the audio data for encoding of plural messages at various levels of distribution (for example, at the production level, the network level and the local affiliate level). Such time division multiplexing of encoded messages substantially restricts bandwidth available for each of the messages and requires a reliable means of determining in each case the permissible time interval for inserting each different message.\nAccordingly, what is needed is a way to encode multiple messages inaudibly in audio data in which one or more such messages are encoded in the audio data at different times and/or levels of distribution which achieves desirably high bandwidth and is easily implemented.\nIt is also desired to provide expanded data communication capability in the limited bandwidth available for ancillary data in an audio channel. It is desired, therefore, to increase the bandwidth afforded by an audio channel to communicate information in the form of ancillary data encoded in the audio data, so that the encoded ancillary data remains inaudible or beneath an acceptable level of audibility when the audio data is reproduced acoustically."}
{"text": "1. Field of the Invention\nThe present invention relates generally to elements that reroute the locations of bond pads on semiconductor devices and, more specifically, to rerouting elements that are configured to be secured to the active surfaces of fabricated semiconductor devices to reroute the bond pad locations thereof. In addition, the present invention relates to methods for designing rerouting elements and to rerouting methods. The present invention also relates to multi-chip modules with semiconductor devices in stacked arrangement and including one or more of the rerouted semiconductor devices, as well as to methods for forming and packaging such assemblies.\n2. Background of Related Art\nIn order to conserve the amount of surface area, or “real estate,” consumed on a carrier substrate, such as a circuit board, by semiconductor devices connected thereto, various types of increased density packages have been developed. Among these various types of packages is the so-called “multi-chip module” (MCM). Some types of multi-chip modules include assemblies of semiconductor devices that are stacked one on top of another. The amount of surface area on a carrier substrate that may be saved by stacking semiconductor devices is readily apparent-a stack of semiconductor devices consumes roughly the same amount of real estate on a carrier substrate as a single, horizontally oriented semiconductor device or semiconductor device package.\nDue to the disparity in processes that are used to form different types of semiconductor devices (e.g., the number and order of various process steps), the incorporation of different types of functionality into a single semiconductor device has proven very difficult to actually reduce to practice. Even in cases where semiconductor devices that carry out multiple functions can be fabricated, multi-chip modules that include semiconductor devices with differing functions (e.g., memory, processing capabilities, etc.) are often much more desirable since the separate semiconductor devices may be fabricated independently and later assembled with one another much more quickly and cost-effectively (e.g., lower production costs due to higher volumes and lower failure rates).\nMulti-chip modules may also contain a number of semiconductor devices that perform the same function, effectively combining the functionality of all of the semiconductor devices thereof into a single package.\nAn example of a conventional, stacked multi-chip module includes a carrier substrate, a first, larger semiconductor device secured to the carrier substrate, and a second, smaller semiconductor device positioned over and secured to the first semiconductor device. Any suitable adhesive may be used to secure the semiconductor devices to one another. The second semiconductor device does not overlie bond pads of the first semiconductor device and, thus, the second semiconductor device does not cover bond wires that electrically connect bond pads of the first semiconductor device to corresponding contacts or terminal pads of the carrier substrate. Such a multi-chip module is disclosed and illustrated in U.S. Pat. No. 6,212,767, issued to Tandy on Apr. 10, 2001 (hereinafter “the '767 Patent”). Due to the use of bond wires to form electrical connections between bond pads and corresponding terminal pads, this type of stacked multi-chip module has been limited to use with semiconductor devices that include peripherally located bond pads.\nU.S. Pat. 5,323,060, issued to Fogal et al. on Jun. 21, 1994 (hereinafter “the '060 Patent”) shows one example where dice of the same size are stacked on top of one another over a circuit board. Bonding wires are connected from the bond pads of each die to corresponding terminal pads on the circuit board. In order to provide clearance for the bond wires that electrically connect bond pads and corresponding terminal pads, however, adjacent semiconductor devices must be spaced apart from one another a significant distance.\nStacked multi-chip modules of other configurations have also been developed. For example, it is known that stacked multi-chip modules may include large semiconductor devices positioned over smaller semiconductor devices and that adjacent semiconductor devices may be staggered relative to one another or have different orientations.\nDifferent electrical connection technologies, including wire bonding, tape-automated bonding (“TAB”), and controlled-collapse chip connection (“C-4”), which results in a so-called flip-chip arrangement are but a few of the ways in which discrete conductive elements may be formed in stacked multi-chip modules. Different electrical connection technologies have also been used in single multi-chip modules, with the bond pads of one semiconductor device being electrically connected to corresponding contact areas of a carrier substrate of the multi-chip module with a different type of discrete conductive element than that used to form electrical connections between the bond pads of another semiconductor device and their corresponding contact areas of the carrier substrate.\nMany semiconductor devices include bond pads that are arranged at central locations on an active surface thereof Examples include semiconductor devices that are configured for use with leads-over-chip (LOC) type lead frames, in which the bond pads are arranged substantially linearly along the centers thereof, as well as semiconductor devices with bond pads disposed in an “I” arrangement. While it may be desirable to use such semiconductor devices in stacked multi-chip modules, the central bond pad placements thereof do not readily facilitate the use of bond wires or other laterally extending discrete conductive elements to electrically connect the bond pads with their corresponding terminal pads of a circuit board that underlies the semiconductor device stack.\nAccordingly, there are needs for apparatus and methods that facilitate the use of semiconductor devices with centrally located bond pads in stacked multi-chip modules. There are also needs for apparatus and methods for reducing the heights of stacked multi-chip modules that include semiconductor devices with peripherally located bond pads."}
{"text": "The present invention relates to a hospital bed. More particularly, the present invention relates to a proning bed which permits rotation of a patient supported on a patient support surface of the bed.\nA bed of the present invention illustratively includes a base, and a support assembly coupled to the base. The support assembly includes first and second support arms located above the base. The apparatus also includes a plurality of latch mechanisms coupled to the first and second support arms, and a patient support surface configured to be coupled to the first and second support arms by the plurality of latch mechanisms. The patient support surface is removable from the first and second support arms to permit transfer of a patient to and from the bed on the patient support surface.\nIn an illustrated embodiment, the support assembly is coupled to a first end of the base. The support assembly includes a rotatable drive mechanism coupled to the first and second support arms for rotating the first and second arms about a longitudinal axis. The first and second support arms may be cantilevered from the support assembly or coupled to a support located at the end of the second base.\nA proning surface is configured to be coupled to the first and second support arms. The proning surface is configured to support the patient in a prone position when the patient support assembly is rotated 180xc2x0 about its longitudinal axis by the drive mechanism. In an illustrated embodiment, a plurality of siderails is coupled to the first and second support arms. The siderails each include a portion which is movable over the patient support surface to form a portion of the proning surface.\nThe illustrated patient support surface includes an outer frame configured to be coupled to the plurality of latch mechanisms to secure the patient support surface to the first and second support arms. The patient support surface also includes a plurality of panels coupled to the outer frame. The panels illustratively include notched portions configured to define handles on the patient support surface. The plurality of panels is pivotably coupled to the outer frame.\nThe illustrated patient support surface further includes at least one hinge to permit articulation of the patient support surface. The patient support surface includes at least one locking member configured to block pivotal movement of the hinge to hold the patient support surface in a generally planar orientation. The support assembly includes an actuator for selectively releasing the locking member to permit articulation of the patient support surface. In an illustrated embodiment, a latching mechanism is configured to engage each locking member. The actuator is configured to move the latching mechanism relative to the first and second support arms to expose the hinge and permit articulation of the patient support surface.\nThe illustrated embodiment of the present invention includes a transfer surface coupled to the base. The transfer surface is movable from a lowered position to an elevated position located adjacent the first and second support arms when the patient support surface is coupled to and removed from the first and second support arms. In one embodiment, the transfer surface is configured to engage a portion of the plurality of latch mechanisms as the transfer surface is moved to the elevated position to open the latch mechanisms for receiving the patient support surface.\nA line management apparatus of the present invention is configured to be coupled to a patient support surface for routing medical lines and hoses. The apparatus includes a body portion having a top edge. The body portion is formed to include a plurality of notches opening along the top edge to receive the lines and hoses and a plurality of apertures located below the notches for receiving additional lines and hoses. A coupler is coupled to the body portion adjacent to the plurality of apertures. The coupler is configured to connect the body portion to the patient support surface\nAdditional features of the invention will become apparent to those skilled in the art upon consideration of the following detailed description of illustrated embodiments exemplifying the best mode of carrying out the invention as presently perceived."}
{"text": "This invention relates to a dual optical fiber device and is particularly concerned with a device for transmitting electromagnetic radiation to a radiation sensitive component and receiving output radiation from said component involving a radiation-transmissible junction of optical fibers which provides a substantially unattenuated output radiation signal. The invention also relates to a method for determining a desired parameter which is a function of said output radiation using said device.\nThe measurement of desired parameters in liquids, particularly in biological systems, is frequently required. For example, the measurement in blood of pH levels and concentration of gases, particularly oxygen and carbon dioxide, is important during surgical procedures, post-operatively, and during hospitalization under intensive care and many devices for the measurement of said physiological parameters have been suggested in the art.\nU.S. Pat. No. 4,003,707, Lubbers et al, and its reissue patent Re 31879, disclose a method and an arrangement for measuring the concentration of gases and the pH value of a sample, e.g. blood, involving the use of a fluorescent indicator at the end of a light-conducting cable which is sealingly covered by or embedded in a selectively permeable diffusion membrane. The radiation transmitted to and emitted from the indicator must be passed through various filtering elements and light elements, including reflectors, beam splitters and amplifiers before any meaningful measurements can be made.\nU.S. Pat. No. 4,041,932, Fostick, discloses a method whereby blood constituents are monitored by measuring the concentration of gases or fluids collected in an enclosed chamber sealingly attached to a skin \"window\" formed by removing the stratum corneum over a small area of the patient's skin. The measurements in the enclosed chamber are made, inter alia, by determining the difference in intensity of light emitted from a fluorescent indicator.\nU.S. Pat. Nos. 4,200.110 and 4,476,870, Peterson et al, disclose the use of a pH sensitive indicator in conjunction with a fiber optic pH probe. In each of these patents the dye indicator is enclosed within a selectively permeable membrane envelope.\nU.S. Pat. No. 4,548,907, Seitz et al, discloses a fluorescent-based optical sensor comprising a membrane immobilized fluorophor secured to one end of a bifurcated fiber optic channel for exposure to the sample to be analyzed.\nMany fluorescent indicators sensitive to pH, and thereby useful for pCO.sub.2 measurements, are known in the art. Examples of useful fluorescent indicators are disclosed in the above patents and also in \"Practical Fluorescence\" by George E. Guilbault, (1973) pages 599-600.\nSensor devices using fluorescent indicators may be used for in vitro or in vivo determinations of components in physiological media. For in vitro determinations the size of the device is normally of no consequence, but for in vivo use, the size of the sensor may be extremely critical and there is an increasing need in the art to miniaturize sensor devices, particularly catheter-type devices, for the in vivo determination of components in physiological media, e.g. blood. However, diminution in size of the components of such devices, particularly in the size of the sensor itself, decreases the strength of the signal emitted by the indicator and consequently presents problems in the detection and measurement of said signal. These problems are aggravated when the detector system requires a multiplicity of components, such as filters, beamsplitters and reflectors to isolate and measure the emitted energy. Each of the said components reduces the emitted signal strength resulting in a sequential loss of measurable signal. Consequently, the more components present in the system, the weaker the final signal strength.\nIt has now been found that the emission signal from radiation-sensitive indicators, particularly fluorescent indicators of the type disclosed in the references discussed above, may be received substantially unattenuated in a suitable detector without the necessity of filters, beam splitters, reflectors or other light elements used in the prior art if the transmitting optical fiber and emission-receiving optical fiber are coupled in a manner according to the present invention, as described hereinafter.\nThe elimination of the need for additional optical elements provided by the improved device of the present invention allows greater miniaturization of sensor systems than that attainable in prior art systems without loss of signal strength."}
{"text": "This disclosure relates to the polishing of magnetic, optical, semiconductor and silicon wafers and more particularly, to polishing compositions and methods for polishing and planarizing silicon wafers.\nThe semiconductor industry uses interconnect metals in forming integrated circuits on semiconductor wafers. These interconnect metals are preferably non-ferrous metals. Suitable examples of such non-ferrous interconnects are aluminum, copper, gold, nickel, and platinum group metals, silver, tungsten and alloys comprising at least one of the foregoing metals. These interconnect metals have a low electrical resistivity. Copper metal interconnects provide excellent conductivity at a low cost. Because copper diffuses into many dielectric materials, such as silicon dioxide or doped versions of silicon dioxide, integrated circuit fabricators typically apply a diffusion barrier layer to prevent the copper diffusion into the dielectric layer. For example, barrier layers for protecting dielectrics include, tantalum, tantalum nitride, silicon nitride, tantalum-silicon nitrides, titanium, titanium nitrides, titanium-silicon nitrides, titanium-titanium nitrides, titanium-tungsten, tungsten, tungsten nitrides and tungsten-silicon nitrides.\nIn the manufacturing of semiconductor wafers, polishing compositions are used to polish semiconductor substrates after the deposition of the metal interconnect layers. Typically, the polishing process uses a “first-step” slurry specifically designed to rapidly remove the excessive interconnect metal. The polishing process then includes a “second-step” slurry to remove the barrier layer. The second-step slurry selectively removes the barrier layer without adversely impacting the physical structure or electrical properties of the interconnect structure. In addition to this, the second step slurry should also display low erosion for dielectrics.\nWhile the “two-step” slurry generally performs satisfactorily in present conditions, it is desirable to have polishing compositions that can be used for planarizing the metal interconnect without substantial dishing of the interconnect metal and scratching of the wafer surface.\nPreparation and evaluation of chemically modified abrasives for chemical mechanical polishing by Partch et al. in an abstract presented at Materials Research Society (MRS) meeting in 2001, discloses coating individual particle cores with a layer of organic or inorganic material. However, Partch et al., disclosed manufacturing the polishing composition with abrasive particles having a large particle size distribution that resulted in scratching and other forms of surface damage.\nU.S. Pat. Pub. No. 2004/0060502 to Singh discloses a plurality of particles that may be coated by weakly or strongly adsorbing surfactants or polymer additives to a core particle. The polymers thus adsorbed onto the abrasive particles can desorb from the particles during the polishing process to alter polishing characteristics and possibly result in damage to the polished surface.\nThere thus remains an unsatisfied demand for polymer-coated particles suitable for chemical mechanical polishing of magnetic, optical, semiconductor or silicon substrates."}
{"text": "This invention relates to devices for the unique marking and identification of various products, e.g., credit cards, machine parts, computer chips, bullets, etc., and to the fabrication and encoding of such devices, and particularly to the secure encoding of devices. By \"secure encoding\" is meant encodings which are difficult to change and/or eradicate, and which survive usage of the encoded product.\nA wide variety of systems have been recently developed for mounting or implanting devices in various products as a means for marking, e.g., identifying the products. Possibly the best known product example is the \"smart card\", comprising a plastic credit card having embedded therein a microelectronic semiconductor device including a programmable memory. Using appropriate writing means, data about the owner of the card, including identification and financial data, can be written into the memory and later read out by insertion of the card into an appropriate read out machine located on a merchant's premises or the like.\nNumerous other examples of systems and devices for uniquely encoding products are known and described, for example, in the following U.S. patents, the subject matter of which are incorporated herein by reference: U.S. Pat. Nos. 5,142,128 (Perkin et al), 4,010,355(Opicella, Jr. et al), 4,839,875 (Kuriyama et al), 4,827,110 (Rossi et al), 5,182,543 (Siegel et al), 5,166,676 (Milheiser), 4,959,515 (Zavrachy et al), 4,685,515 (Anderson et al), 5,175,424 (Lisimaque) and 4,752,776 (Katzenstein). The various devices and systems disclosed in these patents are likely useful, but all contain limitations and disadvantages relating, primarily, to large size, ability to operate in difficult media, strength or ability to withstand acceleration, the need for electrical or other connections, and/or the ability to be embedded in a hidden location in a product.\nFor example, and with reference to certain of the patents in the foregoing list having similarities to the present invention, Zavrachy et al (U.S. Pat. No. 4,959,515) disclose a structure comprising a plurality of cantilevered beams, the free ends of which overlie electrodes on a substrate. Upon the application of electrical charges to the beam ends and electrodes, the beams are deflected, thereby changing the spacing between the beam ends and the electrodes, and thereby providing an electrically detectable event, e.g., a variation in a parameter of an electrical circuit. By selectively disconnecting various ones of the beams from the associated electrical circuits, only selected beams remain to generate electrical signals upon the application of the electrical charges. The pattern of signals thus produced corresponds to a unique code for the device.\nDisadvantages of the device are that it includes numerous electrical circuits all of which have to be individually addressed. The device is thus relatively complex, cannot be hidden in a product, and cannot withstand high temperatures.\nAnderson (U.S. Pat. No. 4,686,515) discloses the use of a strip of magnetostrictive, amorphous metal adapted to be magnetically biased for mechanical resonance at a known frequency. Although not disclosed in the patent, the metal strip or strips are somehow positioned on a marker applied to the product to be identified so that each strip is free to vibrate. By applying to the markers both a dc bias field and an ac interrogation field from a reading mechanism, energy is alternately stored and released with the frequency of the ac field. The quantities of energy storage and release are greatest at the mechanical resonance frequency of the strip, hence an electrical signal is generated having a unique pattern characteristic of the particular resonant frequency of the strip. While the patent refers to the use of \"strips\", the patent is entirely silent as to how such plural strips would be fabricated, encoded, mounted and decoded. Also, the system disclosed is complex and expensive.\nNot included in the foregoing list is U.S. Pat. No. 5,001,933 (Brand), which relates not to an identification system but to a vibration sensor. This patent discloses structure related to the present invention in that cantilevered beams are used, but, as with U.S. Pat. No. 4,959,515, the beams are individually electrically addressed, hence the device and associated systems are comparatively complex and expensive.\nAs explained hereinafter, the present invention provides significant advantages over the prior art."}
{"text": "1. Field of the Invention\nThe present invention is related to development of Web sites and applications. More specifically, the present invention relates to facilitated Internet communications between a client-side and a server-side.\n2. Description of the Related Art\nPrior to Rich Internet Applications, traditional Web applications involved a client-server architecture with all of the processing on the server side and the client-side used to display the HTML web-pages served by the server. Each time a user desired to view a new Web-page, a HTTP request was sent to the server and the requested Web-page was served to the Web browser on the client-side. Such a traditional system is shown in FIG. 1 with a Web-server 1000 on a server side receiving requests over the Internet 1005 from a Web-browser 1003 on a client-side.\nRich Internet Applications, such as Ajax, greatly improved on the traditional client-server architecture by allowing the client machine to dynamically render and partially refresh web pages based on an initial set of instructions from the server, user input, and small amounts of subsequent data dynamically requested from the server. As shown in FIG. 2, the client machine processes Ajax instructions to render a Web page for the user.\nEarly Web applications allowed a user's browser to send a request to a server. The server processed the request and responded to the browser with a Web page. When the user wanted to view a new page, another request was sent to the server and the server responded to the browser with a new Web page. Such a process resulted in a waste of bandwidth since much of the Web contents in the first Web page were also contained in the second web page. The need to resend the same information led to a much slower user interface of a Web application than that of a native application.\nAn emerging technology, called Ajax (Asynchronous and JavaScript XML), was developed for refreshing part of a page instead of refreshing the whole page on every interaction between the user and application. In an Ajax application, when a user submits a form in a page, a script program, usually a JavaScript program, resident on the Web browser receives the user's request and sends a XML (Extended Markup Language) HTTP (Hyper Text Transfer Protocol) request to the Web server in background so as to retrieve only the needed Web contents instead of the whole page and perform corresponding processing to partly refresh the page when receiving a response from the Web server. In this way, the application response time is shortened, because the amount of data exchanged between the Web browser and the Web server is greatly reduced. And the processing time of the Web server is saved because much of the processing is performed at the client side.\nGeneral definitions for terms utilized in the pertinent art are set forth below.\nAjax is the use of dynamic HTML, JavaScript and CSS to create dynamic and usually interactive Web sites and applications. A more detailed explanation of Ajax is set forth in Edmond Woychowsky, AJAX, Creating Web Pages with Asynchronous JavaScript and XML, Prentice Hall, 2007, which is hereby incorporated by reference in its entirety.\nApplets or Java Applets are mini-executable programs named with the .class suffix and are placed on a Web page and provide interactive and multimedia uses.\nApplication Programming Interface (API) is a collection of computer software code, usually a set of class definitions, that can perform a set of related complex tasks, but has a limited set of controls that may be manipulated by other software-code entities. The set of controls is deliberately limited for the sake of clarity and ease of use, so that programmers do not have to work with the detail contained within the given API itself.\nAn Attribute provides additional information about an element, object or file. In a Document Object Model, an attribute, or attribute node, is contained within an element node.\nBehavioral layer is the top layer and is the scripting and programming that adds interactivity and dynamic effects to a site.\nBinding in a general sense is the linking of a library to an application program usually to prevent repetition of frequently utilized code.\nCascading Style Sheets (CSS) is a W3C standard for defining the presentation of Web documents.\nCompiler is a computer program that translates a series of instructions written in one computer language into a resulting output in a different computer language.\nDocument Object Model (DOM) Element is an object contained in a Document Object Model (DOM). The term DOM is generally used to refer to the particular DOM held in the memory region being used by the Web browser. Such a DOM controls the Graphical Respondent Interface (GRI) or Graphical User Interface (GUI). The DOM is generated according to the information that the Web browser reads from the HTML file, and/or from direct JavaScript software instructions. Generally, there exists a unique DOM element for every unique HTML element. DOM elements are sometimes referred to as HTML/DOM elements, because the DOM element exists only because HTML code that was read by the Web browser listed some HTML element that had not previously existed, and thereby caused the Web browser to create that DOM element. Often specific elements of the greater set of HTML/DOM elements are identified by specifying an HTML/DOM checkbox element, or an HTML/DOM text input element. A more detailed explanation of the document object model is set forth in Jeremy Keith, DOM Scripting, Web Design with JavaScript and the Document Object Model, friends of, 2005, which is hereby incorporated by reference in its entirety.\nHyperText Markup Language (HTML) is a method of mixing text and other content with layout and appearance commands in a text file, so that a browser can generate a displayed image from the file.\nHypertext Transfer Protocol (HTTP) is a set of conventions for controlling the transfer of information via the Internet from a Web server computer to a client computer, and also from a client computer to a Web server.\nInternet is the worldwide, decentralized totality of server computers and data-transmission paths which can supply information to a connected and browser-equipped client computer, and can receive and forward information entered from the client computer.\nJavaScript is an object-based programming language. JavaScript is an interpreted language, not a compiled language. JavaScript is generally designed for writing software routines that operate within a client computer on the Internet. Generally, the software routines are downloaded to the client computer at the beginning of the interactive session, if they are not already cached on the client computer. JavaScript is discussed in greater detail below.\nJSON is JavaScript Object Notation format, which is a way of taking data and turning it into valid JavaScript syntax for reconstituting an object at the other end of the transmission protocol.\nMySQL is a relational database management system which relies on SQL for processing data in a database.\nParser is a component of a compiler that analyzes a sequence of tokens to determine its grammatical structure with respect to a given formal grammer. Parsing transforms input text into a data structure, usually a tree, which is suitable for later processing and which captures the implied hierarchy of the input. XML Parsers ensure that an XML document follows the rules of XML markup syntax correctly.\nPlatform is the combination of a client computer, an operating system, and a browser, which together can support Internet access and in particular the operation of interactive forms.\nPresentation layer follows the structural layer, and provides instructions on how the document should look on the screen, sound when read aloud or be formatted when it is printed.\nRendering engine is software used with a Web browser that takes Web content (HTML, XML, image files) and formatting information (CSS, XSL) and displays the formatted content on a screen.\nSerialization places an object in a binary form for transmission across a network such as the Internet and deserialization involves extracting a data structure from a series of bytes.\nSQL (Structured Query Language) is a computer language designed for data retrieval and data management in a database.\nStructural layer of a Web page is the marked up document and foundation on which other layers may be applied.\nUser is a client computer, generally operated by a human being, but in some system contexts running an automated process not under full-time human control.\nWeb-Browser is a complex software program, resident in a client computer, that is capable of loading and displaying text and images and exhibiting behaviors as encoded in HTML (HyperText Markup Language) from the Internet, and also from the client computer's memory. Major browsers include MICROSOFT INTERNET EXPLORER, NETSCAPE, APPLE SAFARI, MOZILLA FIREFOX, and OPERA.\nWeb-Server is a computer able to simultaneously manage many Internet information-exchange processes at the same time. Normally, server computers are more powerful than client computers, and are administratively and/or geographically centralized. An interactive-form information-collection process generally is controlled from a server computer, to which the sponsor of the process has access.\nWorld Wide Web Consortium (W3C) is an unofficial standards body which creates and oversees the development of web technologies and the application of those technologies.\nXHTML (Extensible Hypertext Markup Language) is a language for describing the content of hypertext documents intended to be viewed or read in a browser.\nXML (Extensible Markup Language) is a W3C standard for text document markup, and it is not a language but a set of rules for creating other markup languages.\nThere are three types of JavaScript: 1) Client-side JavaScript; 2) Server-side JavaScript; and 3) Core JavaScript. Client-side JavaScript is generally an extended version of JavaScript that enables the enhancement and manipulation of web pages and client browsers. Server-side JavaScript is an extended version of JavaScript that enables back-end access to databases, file systems, and servers. Core JavaScript is the base JavaScript.\nCore JavaScript includes the following objects: array, date, math, number and string. Client-side JavaScript and Server-side JavaScript have additional objects and functions that are specific to client-side or server-side functionality. Generally, any JavaScript libraries (.js files) created in core JavaScript can be used on both the client and the server without changes. Client-side JavaScript is composed of a Core JavaScript and additional objects such as: document, form, frame and window. The objects in Client-side JavaScript enable manipulation of HTML documents (checking form fields, submitting forms, creating dynamic pages) and the browser (directing the browser to load other HTML pages, display messages). Server-side JavaScript is composed of Core JavaScript and additional objects and functions for accessing databases and file systems, and sending email. Server-side JavaScript enables Web developers to efficiently create database-driven web applications. Server-side JavaScript is generally used to create and customize server-based applications by scripting the interaction between objects. Client-side JavaScript may be served by any server but only displayed by JavaScript-enabled browsers. Server-side JavaScript must be served by a JavaScript-enabled server but can be displayed by any browser.\nUnited States Patent Application Publication Number 20010037359 describes a system and method for a server-side browser including markup language graphical user interface, dynamic markup language rewriter engine and profile engine. The system includes a user computer and a destination server computer separated by a server computer hosting a server-side browser (SSB). The SSB includes a markup language graphical user interface (MLGUI), a dynamic markup language rewriter engine (DMLRE) and a profiling engine (PE). The SSB may be configured as an intermediary infrastructure residing on the Internet providing customized information gathering for a user. The components of the SSB allow for controlling, brokering and distributing information more perfectly by controlling both browser functionality (on the client-side) and server functionality (on the destination site side) within a single point and without the necessity of incremental consents or integration of either side.\nHowever, current technologies that operate Server-side JavaScript fail to offer complete interactions which are the hallmark of rich web sites and applications. When writing software, it is often convenient to group the code according to the area of functionality it provides. But the group might span code that needs to run on the server, code that needs to run on the client, code that needs to run on both, or run on the server and get stored there and have a proxy for it created on the client, etc."}
{"text": "1. Field of the Invention\nThis invention relates generally to folding, collapsible structures, and more particularly relates to an awning assembly for a shelter having a canopy portion with an upper framework, and at least two adjacent legs supporting the canopy portion, the awning assembly adapted to be removably mounted to at least two adjacent legs.\n2. General Background and State-of-the-Art\nTemporary shelters that can be easily transported and rapidly set up at emergency sites can be particularly useful in providing temporary care and housing. Such shelters can also be useful for non-emergency outdoor gatherings, such as for temporary military posts, field trips, and the like. It would be desirable to provide an improved collapsible shelter with a multipurpose awning that can be moved to different positions to provide shade and to serve as a holder for signage. It would also be desirable to provide a modular multipurpose awning that is adapted to mounted to an existing shelter and that can be moved to different positions to provide shade and to serve as a holder for signage. The present invention fulfills these and other needs."}
{"text": "In the electronics industry, a continuing need is to further and further reduce the size and weight of electronic devices while simultaneously increasing performance and speed. Cellular telephones, personal data devices, notebook computers, camcorders, and digital cameras are but a few of the consumer products that require and benefit from this ongoing miniaturization of sophisticated electronics.\nIntegrated circuit (“IC”) assemblies for such complex electronic systems typically have a large number of interconnected IC chips. The IC chips, commonly called dies, are usually made from a semiconductor material such as silicon or gallium arsenide. Photolithographic techniques are used to form the various semiconductor devices in multiple layers on the dies.\nDies are encapsulated in a molded plastic package that has connectors or leads on the exterior of the package that function as input/output terminals for the die inside the package. The package includes a leadframe, a die mounted on the leadframe, and wires.\nThe die is conventionally mounted to the top surface of the leadframe with, for example, a layer of an adhesive or an adhesive film, and then electrically connected to the leadframe by a number of fine, conductive wires, typically gold (Au) or aluminum (Al), that electrically connect the die to the leadframe. The wires are attached to the die at the bonding pads of the die, which are located around the periphery of the die.\nAfter the wires are attached, the die, the leadframe, and the conductive wires are encapsulated in a mold compound, such as plastic or epoxy. The encapsulation protects the leadframe, the die, and the fine conductive wires from physical, electrical, moisture, and/or chemical damage.\nThe encapsulation process begins by placing the leadframe, the die, and the fine conductive wires in a mold. Next, a mold compound is injected into the mold. The mold compound flows through the mold, encasing the leadframe, the die, and the conductive wires.\nTypically, a mold encapsulates multiple semiconductor devices at the same time. A two part mold mounted on a hydraulic press is generally used.\nInitially the two halves of the mold are held apart. One or more lead frames containing semiconductor devices are placed in an open half of the mold. The hydraulic press is then actuated and the mold closed, forming a cavity around each semiconductor die. It is not unusual for a mold to contain thousands of cavities. Each of these cavities is connected by one or more gates, channels, and runners to one or more softened plastic central reservoirs or pots. A worm screw or ram compresses the plastic so that it flows into the cavities. As soon as the plastic has hardened, the mold is opened and the molded plastic packages removed.\nFrequently, moveable pins are built into the mold to align the mold parts, to hold the lead frames in a particular location during molding or to provide automatic ejection of the encapsulated parts. Also, the mold may contain other moving parts such as variable gates, vents, and dams. Thus, molds for encapsulating electronic parts, particularly semiconductor parts, are often very complicated.\nIn order to push liquefied plastic from the reservoirs into the many cavities, it is frequently necessary to inject or transfer the plastic at very high pressures. If the mold halves fail to seal tightly against each other or against the lead frames, undesired or unintended crevices may be present therebetween. At such high pressures, the crevices fill with plastic during encapsulation, producing unwanted thin webs of plastic. These thin webs of plastic are referred to as “flash” and can result from such imperfect sealing of the mold. Before the encapsulated electronic devices can be used, this flash must be removed. This increases the cost of manufacture and is undesirable. Also, flash is a significant cause of mold wear, requires additional labor for mold cleaning between molding cycles, and increases mold down-time.\nIn order to minimize flash, great pains are generally taken to machine the mating surfaces of the mold halves flat and parallel where they are to seal. Usually, they are carefully inspected for planarity during manufacture and after installation in the press. Powerful hydraulic cylinders in the mold force the mold halves tightly against the lead frames and each other. However, the force that can be applied in an effort to seal the mold is limited, since excessive force causes coining of the lead frames and rapid mold wear.\nAs a result, adhesives, such as adhesive tape, are commonly used to secure leadframes and reduce mold flash. However, mold flash continues to occur. Furthermore, adhesives must be cleaned from leadframes, increasing process steps and cost. Despite these efforts, flash continues to occur, even in the most carefully fabricated molds. There is therefore also a need to selectively shield integrated circuits compatibly with reducing mold flash.\nAfter encapsulation, integrated circuits are arranged on a printed circuit board with other integrated circuits and electronic components. Small electronic devices require the integrated circuits to be very close together. However, some integrated circuits are sources of electromagnetic interference and must be spaced further away from other integrated circuits. This required spacing increases the size of the small electronic device being assembled. What is needed is a way to manufacture shielded integrated circuits while reducing mold flash.\nThus, a need still remains for improved encapsulation methods for shielded and unshielded integrated circuits that reduce size, reduce weight, and eliminate unwanted flash.\nSolutions to these problems have been long sought but prior developments have not taught or suggested any solutions and, thus, solutions to these problems have long eluded those skilled in the art."}
{"text": "In the latter half of the twentieth century, there began a phenomenon known as the information revolution. While the information revolution is a historical development broader in scope than any one event or machine, no single device has come to represent the information revolution more than the digital electronic computer. The development of computer systems has surely been a revolution. Each year, computer systems grow faster, store more data, and provide more applications to their users.\nA modern computer system typically comprises one or more central processing units (CPU) and supporting hardware necessary to store, retrieve and transfer information, such as communication buses and memory. It also includes hardware necessary to communicate with the outside world, such as input/output controllers or storage controllers, and devices attached thereto such as keyboards, monitors, tape drives, disk drives, communication lines coupled to a network, etc. The CPU or CPUs are the heart of the system. They execute the instructions which comprise a computer program and directs the operation of the other system components.\nFrom the standpoint of the computer's hardware, most systems operate in fundamentally the same manner. Processors are capable of performing a limited set of very simple operations, such as arithmetic, logical comparisons, and movement of data from one location to another. But each operation is performed very quickly. Sophisticated software at multiple levels directs a computer to perform massive numbers of these simple operations, enabling the computer to perform complex tasks. What is perceived by the user as a new or improved capability of a computer system is made possible by performing essentially the same set of very simple operations, but using software having enhanced function, along with faster hardware.\nIn the very early history of the digital computer, computer programs which instructed the computer to perform some task were written in a form directly executable by the computer's processor. Such programs were very difficult for a human to write, understand and maintain, even when performing relatively simple tasks. As the number and complexity of such programs grew, this method became clearly unworkable. As a result, alternate forms of creating and executing computer software were developed. In particular, a large and varied set of high-level languages was developed for supporting the creation of computer software.\nHigh-level languages vary in their characteristics, but all such languages are intended to make it easier for a human to write a program to perform some task. Typically, high-level languages represent instructions, fixed values, variables, and other constructs in a manner readily understandable to the human programmer rather than the computer. Such programs are not directly executable by the computer's processor. In order to run on the computer, the programs must first be transformed into a form that the processor can execute.\nTransforming a high-level language program into executable form requires that the human-readable program form (source code) be converted to a processor-executable form (object code). This transformation process generally results in some loss of efficiency from the standpoint of computer resource utilization. Computers are viewed as cheap resources in comparison to their human programmers. High-level languages are generally intended to make it easier for humans to write programming code, and not necessarily to improve the efficiency of the object code from the computer's standpoint. The way in which data and processes are conveniently represented in high-level languages does not necessarily correspond to the most efficient use of computer resources, but this drawback is often deemed acceptable in order to improve the performance of human programmers.\nWhile certain inefficiencies involved in the use of high-level languages may be unavoidable, it is nevertheless desirable to develop techniques for reducing inefficiencies where practical. This has led to the use of compilers and so-called “optimizing” compilers. A compiler transforms source code to object code by looking at a stream of instructions, and attempting to use the available resources of the executing computer in the most efficient manner. For example, the compiler allocates the use of a limited number of registers in the processor based on an analysis of the instruction stream as a whole, and thus hopefully minimizes the number of load and store operations. An optimizing compiler might make even more sophisticated decisions about how a program should be encoded in object code. For example, it might determine whether to encode a called procedure in the source code as a set of in-line instructions in the object code.\nEven with all the compilation and associated high-level language tools available to the programmer, there are still some types of executable programming code, typically low-level operating system kernel functions, which are of such critical importance that they are manually programmed at a much lower level to achieve greater computer resource efficiency. At these lower levels, the programmer may decide how to represent data, allocate registers, assign storage addresses, and do other tasks often performed by the compiler or optimizing compiler.\nA typical program contains many places at which flow of execution may diverge or converge, and many potential paths in the flow of program execution exist. For a typical program, many of these paths are rarely if ever used, while a relatively small number of the paths are utilized frequently. Rarely used paths may exist to handle special cases or errors, or may be unintentional side effects of the way in which a program was written. A program will generally perform more efficiently if the bulk of the system's resources are allocated to the most frequently used paths. For example, variables which occur in the most frequently used paths should be given preferences in the allocation of registers over variables which occur in the rarely used paths. Unfortunately, it is difficult for a compiler or optimizing compiler to know in advance which are the frequently used paths, since whether a path is frequently used or otherwise depends on the input data. One of the reasons that programming code written by a programmer at a low level tends to outperform code which is written at a higher level and compiled to object form is that the programmer usually knows better than the compiler which paths will be most frequently used.\nGenerally, it is possible to produce more efficient object code, and particularly to produce more efficient object code using an optimizing compiler, if it can be known in advance what the pattern of usage of the various code paths will be.\nIt is possible to collect data from actual or simulated run-time execution of a computer program in order to determine experimentally the frequency of execution of the various paths of a program. Such data is referred to herein as program execution profile data, or simply profile data for short.\nCommonly, collection of profile data is accomplished by inserting special instructions into the program to collect data at key points. These instructions are referred to herein as “instrumentation instructions”, or “hooks”. A hook, which could be a single instruction or a set of instructions (including a called procedure) causes some record to be made each time it is encountered during execution of a program. Typically, the hook causes a corresponding counter to be incremented, although a record could take some other form.\nA complete and accurate picture of the performance of a computer program requires that the frequency of taking each possible path in the flow of control be known. Because a typical computer program contains a very large number of possible paths, placing instrumentation hooks in every such path to measure flow is a significant burden. However, it is not necessary to directly measure every path. Mathematical techniques exist for determining a subset of the possible paths for instrumentation, from which the frequency of execution of the remaining unmeasured paths can be inferred. These techniques involve the construction of a control glow graph (CFG), which is a directed graph in which each node represents a basic block of code (i.e., a set of sequential instructions having only one entry point and no branches except at the end) and each arc represents a possible path for transfer of control from one block to another (by branching or by fall-through). The frequency of taking a path (arc) in the control flow graph is represented as an arc weight. It is assumed that flow in the graph is conserved, i.e., the sum of the arc weights of all arcs entering any node is equal to the sum of the arc weights of the arcs leaving the node. From a control flow graph, a spanning tree of arcs can be determined, such that the arc weight of any arc can be inferred from the weights of the arcs that are not in the spanning tree, based on the assumption that flow is conserved. Therefore, if instrumentation hooks are inserted only in the paths represented by arcs not in the spanning tree, the frequency of taking other paths can be inferred. Typically, a spanning tree can be constructed such that only 30%-40% of the arcs in the control flow graph need be instrumented, thus realizing a considerable reduction in the number of instrumentation hooks required.\nHowever, even 30%-40% of the possible paths in a program often represents a very large number of paths. To minimize the deleterious effect of instrumentation hooks on program performance, instrumentation code should be as simple as possible. Specifically, a given hook usually increments a single counter in memory only, without performing other operations. A separate counter is associated with each hook. Counter values are examined and used to derive additional data only after data collection from the program ceases.\nWhere multiple processes execute the same instrumented program code simultaneously, the simplicity of the instrumentation code can lead to errors. The multiple processes need to access and increment the same counters, yet the instrumentation code has no protection against contention. If two processes both attempt to read, increment, and write back to the same counter simultaneously, one of the increments may be lost. This effect is referred to as “counter contention”.\nIf all possible paths (control flow arcs) in a program are instrumented, the effects of counter contention is typically small. However, as explained above, instrumenting all paths is very burdensome. Where the arc weights of many paths are inferred from a smaller number of measured paths, errors in the measured paths due to counter contention can be propagated a significant distance in the graph. This may cause counter errors to propagate into code paths which are infrequently or never taken. A compiler attempting to optimize code based on such data may skew the optimization in favor of such paths, to the detriment of other areas of the programming code.\nA need exists for a method and apparatus for obtaining more accurate profile data, without the burden of overly complex instrumentation code or larger numbers of instrumentation hooks."}
{"text": "The present invention relates to personal care products, such as wipes and absorbent articles, that are capable of providing a skin health benefit to the user. More particularly, the present invention relates to wet wipes comprising at least one botanical compound that selectively controls the growth and/or adherence of flora on the skin. The botanical compounds described herein and suitable for incorporation into a wipe or absorbent article may selectively promote the adherence of healthy flora to the surface of skin or mucosa, selectively inhibit the growth of problem flora on or around the skin surface, selectively inhibit the adherence of problem flora to the surface of skin or mucosa, or provide multiple combinations of these effects.\nA variety of flora, both beneficial and pathogenic, may be found on the skin at any given time. Problem flora, such as pathogenic bacteria and yeast, have been associated with numerous ailments, including skin infections, diaper rash, urinary or vaginal infections, and malodors, and are associated with various irritants, such as proteases, lipases, carbohydrases, lipopolysaccharides (“LPS”), volatile organic compounds (“VOCs”), and other bio-molecules. In contrast, healthy bacteria attach to the skin, and may provide a variety of benefits, including playing a role in preventing pathogenic organisms from colonizing or growing.\nThe disruption in the protection of healthy flora may result in the colonization of other organisms, which in turn may cause irritation, infection, and diseases on or near the skin surface. One example of a healthy bacterium is Lactobacillus acidophilus, which colonizes the vaginal epithelium during child-bearing years, and inhibits the growth of pathogens. Another skin bacterium which can have positive health attributes is Staphylococcus epidermidis. S. epidermidis is a normal microbial inhabitant of human skin. In most cases, strains of S. epidermidis are nonpathogenic and play a protective role in their host. For instance, S. epidermidis appears to prevent colonization of the skin by dermatophytic fungi, and may alter the production of irritating metabolites. In limited circumstances, S. epidermidis, however, can become an opportunistic pathogen by spreading into the blood through breaks in skin barriers. Individuals most susceptible to S. epidermidis infection are intravenous drug users, newborns, elderly, and those using catheters or other artificial appliances. Regardless, S. epidermidis is generally regarded as a bacteria that provides a health benefit.\nIn contrast, an example of a common pathogenic flora is Candida albicans. When in yeast form, C. albicans naturally inhabits the human digestive tract, and is a normal part of bowel flora. However, under certain conditions, such as a weakened immune system, a high sugar diet, or improper pH in the digestive system, among others, C. albicans may shift from the yeast stage, to an invasive mycelial fungal form. The mycelial fungal form of C. albicans binds readily to the skin and mucus membranes, and is associated with a wide variety of infections, including vaginal yeast infections, oral infections, and diaper rash.\nOther common pathogenic flora include Pseudomonas aeruginosa, Proteus mirabilis, and Escherichia coli. Like C. albicans, E. coli naturally inhabit the human gastro-intestinal tract. However, E. coli can be pathogenic, and are responsible for several types of infections in humans, including urinary tract infections. For example, E. coli can colonize from the feces or perineal region, and ascend the urinary tract to the bladder, causing infection and irritation.\nP. aeruginosa is an opportunistic pathogen of humans, and is resistant to many antibiotics. P. aeruginosa may infect almost any type of compromised tissue, and can cause a variety of infections, including urinary tract infections, respiratory system infections, dermatitis, soft tissue infections, and bacteremia, as well as a variety of systemic infections, particularly in immunosuppressed patients.\nP. mirabilis is a Gram negative rod that is widely distributed in nature and easily isolated in the feces of most animals, but is hardly ever found in high numbers unless the normal intestinal microflora is altered. P. mirabilis can cause invasive diarrhea and severe infections of the upper urinary tract, and has also been implicated in infections of blood and wounds.\nTo date, skin cleaning products, such as wet wipes and dry wipes, have primarily cleaned the skin by attempting to remove and/or kill all flora present on the skin, regardless of whether the flora are potentially beneficial or potentially harmful. For example, numerous commercially available wet wipes comprise at least one antimicrobial compound, such as an organic acid, which is typically used in combination with a surfactant to kill flora on the skin surface. Although typically effective in killing flora located on the skin surface, the antimicrobial wipe kills all flora and does not, and cannot, distinguish between “good” and “bad” flora.\nBased on the foregoing, it would be desirable to provide a skin health benefit for inhibiting the growth or adherence to the skin of problem flora, while maintaining, or even enhancing the adherence of healthy flora. Products, such as wipes or absorbent articles, that comprise one or more compounds capable of selectively controlling the growth and/or adherence of flora to the skin would thus be desirable, and could impart a beneficial health effect to the user."}
{"text": "The removal of dirt and debris from streets, parking lots, airport runways, factory floors, and other similar paved surfaces, through the use of various types of street sweeping vehicles or factory sweeping vehicles, as may be the case, has been known for many years. As is well known in the industry, street sweeping vehicles, also known as mechanical street sweepers, employ a sweeping broom to remove dirt and debris from a surface and a conveying type elevating mechanism to lift the debris several feet up and deposit it into a hopper. Also, as is well known in the industry, factory sweeping vehicles, also known as factory sweepers, employ a sweeping broom to remove dirt and debris from a surface and sweep the debris several feet up and deposit it into a hopper. For the sake of brevity, clarity and simplicity, such vehicles will be generally referred to in this document as surface sweeping vehicles.\nIn surface sweeping vehicles, a pair of counter-rotating brushes sweep dirt and debris inwardly to underneath the central area of the sweeper and an elongate cylindrically-shaped sweeping broom that rotates about a horizontal axis sweeps the dirt and debris forward up and onto a conveyor. The conveyor deposits the dirt and debris into a hopper for subsequent controlled dumping from the hopper. Such mechanical broom sweepers can remove large amounts of dirt and debris from a paved surface quite quickly and can generally remove large pieces of debris quite readily. However, they cannot contain fine particulate matter that has become airborne, without the use of water for dust suppression. The use of water is undesirable as it creates two problems. A covering of wet dirt remains on the surface behind the surface sweeping vehicle. During the warm months, when the water in this wet dirt evaporates, significant amounts of dried small particulate matter from the wet dirt become air borne. Also, water cannot be used in cold winter months because the water tends to freeze on the surface being swept, thus creating unsafe conditions, and tends to freeze in tank, lines and water pipes.\nIt is widely accepted in the industry that the containment of dust generated during the street cleaning operation is extremely difficult, especially the containment of dust having a particle size under ten microns, without using water. Virtually all street sweepers--that is to say mechanical street sweeping vehicles employing sweeping brooms as the primary means for removing dirt and debris from a road surface--have an inherent problem with containment of dust, especially dust having a particle size under ten microns. With street sweeping vehicles, it is common to use water for dust suppression. However, water can be used only during warm weather when the water will not freeze, and thus, such sweepers are often avoided altogether in many colder climates. Moreover, after a wet road surface has been swept by a sweeping broom, the wet road surface dries and leaves behind a residual fine dust that ultimately becomes airborne, commonly in the form of very fine dust having a particle size under ten microns, which is highly unacceptable.\nRecently, it has become increasingly important for environmental reasons to not just fully remove dirt and debris during a street cleaning operation, but to remove dust and other particulate matter, especially particles less than about ten microns. In many jurisdictions, there are strict environmental laws pertaining to the removal and containment, during a street cleaning operation, of particulate matter having a size of less than ten microns, which is essentially breathable particulate.\nSome mechanical type street sweeping vehicles employ a rotating broom that throws dirt and debris into a \"squeegee\" type elevator for deposit into a hopper. The elevator overthrows the dirt and debris into a hopper behind the elevator. The filtration system is located within the hopper, directly over the entire volume of the hopper. Air from the hopper is drawn upwardly through the filters, thus always creating a negative pressure within the hopper, elevator, and sweeping broom areas, and is expelled through a small additional filter to the atmosphere. During use, the filters trap dust and fine debris from the air stream created by the fan. A vibrating mechanism shakes the filters to shake loose the trapped dust and debris and deposit it into the hopper. However, since the air stream is continuously flowing at an aggressive pace through the filters, the trapped dust and debris can tend to clog the filters. In order to properly clear the filters completely, the sweeping operation must be stopped, and the filters vibrated until they are clear. Often, additional cleaning and maintenance of the filters may be necessary.\nMoreover, due to the inherent positioning of the filter in the hopper, such filters are susceptible to moisture carryover, which can lead to clogging of the filters, thereby possibly causing cessation of dust control.\nThe only mechanical street sweeping vehicles that can employ an auxiliary debris suction hose to facilitate the supplemental suctioning of debris, must first stop their street sweeping operation due to their inherent design and operating characteristics.\nIt is an object of the present invention to provide a dust controlling apparatus for use in a mechanical surface sweeping vehicle, which dust controlling apparatus substantially precludes dust that is generated during a street cleaning operation from being expelled to the atmosphere.\nIt is an object of the present invention to provide a dust controlling apparatus for use in a mechanical surface sweeping vehicle, which dust controlling apparatus substantially precludes dust that is generated during a street cleaning operation from being expelled to the atmosphere, without the use of water for dust suppression.\nIt is another object of the present invention to provide a dust controlling apparatus for use in a mechanical surface sweeping vehicle, which dust controlling apparatus substantially precludes dust that is generated during a street cleaning operation, and has a particle size of less than ten microns, from being expelled to the atmosphere.\nIt is a further object of the present invention to provide a dust controlling apparatus for use in a mechanical surface sweeping vehicle, which dust controlling apparatus substantially precludes dust that is generated during a street cleaning operation from being expelled to the atmosphere, during the cleaning in both wet and dry street conditions.\nIt is yet another object of the present invention to provide a surface sweeping vehicle that can suction dust and various forms of debris while sweeping a surface with a sweeping broom.\nIt is yet another object of the present invention to provide a surface sweeping vehicle having a dust and debris separator that does not become clogged with moist debris."}
{"text": "This invention relates to an indicator for showing the discharged voltage of a stroboscope used with a camera.\nA photographic stroboscope is generally designed for the discharge of electric energy stored in a discharge capacitor through a flash tube. At the time of discharge, the flash tube sends forth a flash. The intensity of the flash varies with the static energy or voltage stored in the discharge capacitor. To ensure, therefore, a flash having a desired intensity, it is necessary to examine the voltage of the discharge capacitor. In other words, an indicator is needed which shows the voltage of the discharge capacitor.\nA known indicator designed to meet this requirement is the type in which a neon discharge tube is connected to the discharge capacitor. The neon discharge tube requiring a certain discharge space is subject to limitations in being rendered compact. Accordingly, considerable difficulties are encountered in fitting a neon discharge tube type voltage indicator for stroboscope into a camera finder. In contrast, an indicator using a light-emitting diode (LED) in place of a neon discharge tube can be built in a camera finder, enabling a camera as a whole to be easily handled, because the LED can be made very compact. However, a larger amount of current is necessary for the LED to be lighted than when the neon discharge tube is lighted. Discharge current required for the neon discharge tube is about 0.01 mA. Whereas drive current of about 1 mA is required for the LED to give off a light. Therefore, direct replacement of the neon discharge tube used with the prior art indicator by a LED results in a noticeable increase in power consumption, namely, a decrease in the life of a dry cell used as a power source of a stroboscope."}
{"text": "The present invention relates to a valve-moving control apparatus for intermitting operation of a predetermined number of valves of cylinders of an engine in accordance with a load applied to the internal combustion engine and a method therefor.\nIn an internal combustion engine for intermitting operation of a predetermined number of valves of cylinders in accordance with the load applied to the engine, when an opening degree of a throttle is low at a certain rotational speed, a large amount of axis torque output is obtained by intermitting actuations of some valves of the cylinders and when the open degree of the throttle is high at this rotational speed, a large amount torque output is obtained by actuating all valves.\nAccordingly, a switch for intermitting actuations of some valves of the cylinders or actuating all valves is usually operated at a timing (cross point) when the torque output of an internal combustion engine in which some valves are intermitted, and the torque output of the engine in which all valves are actuated, are equal to each other at the same opening degree of a throttle as shown in FIG. 9. When the switch is shifted at the cross point, torque shock does not occur, since the torque output does not change.\nHowever, if an internal combustion engine in which some valves of the cylinders are intermitted is changed to actuate all the valves, inner pressure at the intake pipe of all the cylinder is very high immediately after the operation is changed, and a large volume of air is suddenly intaken into all the cylinders, so that a large torque occurs and its shock is transmitted to an internal combustion engine's mounts, a driving mechanism and wheels of the vehicle in which the engine is mounted.\nThe shock caused by such a torque difference is not very much when a gear is shifted to high speed traveling such as the third, the fourth gear and the fifth gear, for which a gear change ratio is relatively low. On the other hand, the shock is largely transmitted to a chassis when the gear is shifted to a low speed traveling, such as the first or the second gear, for which the gear change ratio is large."}
{"text": "The present invention relates to a digital-signal-processing circuit and a display apparatus and a liquid-crystal projector that employ the digital-signal-processing circuit. More particularly, the present invention relates to a digital-signal-processing circuit carrying out gamma correction on a digital video signal for driving a display device exhibiting a non-linear optical response characteristic and to a display apparatus and a liquid-crystal projector that employ the digital-signal-processing circuit.\nA display apparatus exhibits a device-peculiar non-linear optical response to an input voltage. An example of such a display device is a liquid-crystal display apparatus employing liquid-crystal cells distributed among pixels to serve as electro-optical devices. An example of the optical-response characteristic of the display apparatus is a characteristic representing a relation between the transmittance and the applied voltage, as shown in FIG. 5. Such a characteristic is referred to as a V-T characteristic.\nFrom a gradation-recognition characteristic of a human being, on the other hand, a desirable characteristic of a picture display apparatus display luminance, such as the transmittance, is an exponential function with respect to the level of an input voltage, as shown in FIG. 6. In order to obtain the desirable characteristic shown in FIG. 6 from a picture display apparatus with the V-T characteristic shown in FIG. 5, the voltage to be applied to the liquid-crystal device must be generated appropriately from the level of an input signal, as indicated by the non-linear relation shown in FIG. 7. The non-linear relation shown in FIG. 7 is referred to as a gamma-correction curve for an input digital video signal.\nGenerally known examples of the gamma correction include multi-break-point correction and correction based on an LUT (Look-up Table). These corrections entail the use of an analog or digital circuit. The LUT-based correction using a digital circuit has the problem of a large circuit scale. With a higher degree of IC integration achieved in the recent years, however, the limitation imposed by the large scale of the circuit is reduced. In addition, offering the merit of a high degree of precision, the LUT-based correction using a digital circuit has been becoming most popular.\nIn accordance with a related-art technology for implementing the LUT-based digital gamma correction, a digital gamma-correction circuit is used for carrying out the gamma correction as well as the transformation N≧n+2, where notation n denotes the bit count of the input digital video signal supplied to the digital gamma-correction circuit and notation N denotes the bit count of the digital video signal output by the digital gamma-correction circuit. For details of the digital gamma-correction circuit, refer to documents such as Japanese Patent Laid-open No. 2000-20037. This related-art technology is intended for avoiding gradation deterioration in a gray zone.\nTo put it in detail, the gamma-correction curve representing a relation between the voltage to be applied to the liquid-crystal device and the level of an input signal has large gradients in a black-side zone and smaller ones in the gray zone, as shown in FIG. 7. These gradients indicate that, with the output-bit count set at the same value as the input-bit count, the gradation in the gray zone cannot be kept at the same degree of precision as the input gradation. The related-art technology described above carries out the gamma correction as well as the transformation N≧n+2 in order to avoid gradation deterioration in the gray zone, where notation n denotes the number of bits input to the LUT and notation N denotes the number of bits output from the LUT.\nWith the related-art technology described above, however, the number of output bits rises with an increase in input-bit count. There is thus raised the problem that the output-pin count of a signal-processing IC for carrying out the gamma correction also increases. In addition, the input-pin count of a D/A converter provided at a later stage rises as well. As a result, the circuit scales of the ICs also increase and the power consumption rises as well, raising the problem of an increased amount of unnecessary radiation. Furthermore, the related-art technology described above is no more than a technology for avoiding gradation deterioration in a gray zone where the gradient of the gamma correction curve is small."}
{"text": "Human alpha-thrombin appears to have growth-promoting activity for a wide variety of cells from various tissues. For example, alpha-thrombin has been shown to initiate proliferation of fibroblastic cells in culture without addition of serum or other purified growth factors, to synergize with epidermal growth factor in certain hamster fibroblasts and human endothelial cells, to initiate cell division or DNA synthesis in mammalian lens epithelial and spleen cells and actuate monocytes and neutrophils. Yet, the use of thrombin as a growth factor and its potential importance to wound healing has not been widely acclaimed. In part, this may be due to the complexity of thrombin's involvement with coagulation, platelet activation, and initiation of cell proliferation as well as to the complex regulation of thrombin and thrombin-like molecules by serum protease inhibitors and by cell-released protease nexins. This complexity and high degree of physiologic regulation, however, supports the potential importance of this initiation pathway in wound healing.\nThrombin may also play a role in both normal revascularization and migration of cells from the blood to the site of injury and the abnormal metastasis and angiogenesis associated with tumors. The ability of thrombin to increase endothelial cell proliferation and alter the barrier function of blood vessels may contribute to angiogenesis and inflammation at sites of tissue injury.\nThrombin derivative peptides have been described by the present inventors for the agonizing and antagonizing thrombin and/or thrombin receptor activity, such as in the treatment of wounds. U.S. Pat. No. 5,500,412 or 5,352,664, the contents of which are incorporated herein by reference in their entirety. However, the patent does not teach the novel use of the thrombin derivative peptides for the treatment of damaged cardiac tissue, for revascularization, or for inhibition of vascular occlusion and restenosis."}
{"text": "Control circuits constructed with electronic components, and particularly control circuits having fault protection, have for some time been a part of the state of the art and are used, for example, for controlling presses, punching presses, automated systems and the like. These control circuits pick up fault conditions by using self-monitoring logic circuits, for example, two-channel logic circuits with antivalence monitoring (see German Offenlegungsschrift 27 21 270), or by cyclic monitoring of safety-related components (see German Offenlegungsschrift 26 33 322 or British Pat. No. 1,527,857). Such control circuits are designed for monitoring certain operational conditions such as single-stroke operation with eccentric presses, and in this they guarantee a high degree of operational reliability.\nThe actuators ad/or the power amplifiers preceding the presses, which actuators and/or power amplifiers are driven by the control circuits, are predominantly provided with semiconductor components. These semiconductor components can be connected to a load voltage supply circuit in such a manner that, if the control circuits and the machines to be controlled are functioning properly, the actuators which cause movements will be driven.\nIn the case of a fault, however, as a rule the load volatage supply circuits are cut off. This cut-off is a result of action by the control circuit which is protected against faults. Under special operating conditions, for example, during setting-up, the control circuit can be used in such a manner that it is possible to operate the machine even though certain safety circuits (e,g, protective screen up) are deactivated. During such an operating condition the semiconductor components which switch the load voltage supply circuit are connected direactly into the active circuit. Therefore, a failure of a semiconductor component can lead to an unplanned movement, for example, in an automated system, which greatly endangers the operator during setting-up operations."}
{"text": "This invention pertains to the field of digital optical communications.\nA typical digital optical communications system includes a transmitter, an optical channel (e.g., optical fiber), and a receiver. Optionally, the optical channel may have one or more elements such as amplifiers to compensate for attenuation due to long distances and other dispersive effects. The transmitter is often constructed of a light source generating an optical carrier and an electro-optical modulator. A common electro-optical modulator is the Mach Zehnder (M-Z) modulator. The M-Z modulator operates by dividing incoming light into two optical paths. Into one or both paths is placed a material that, in response to a control voltage, alters the phase of the light traveling through it. The two paths are then optically recombined to generate the final output. Phase differences between the two paths cause destructive interference, operating to reduce the intensity of the output. When the two paths are exactly one-half wavelength out of phase (a radians), destructive interference is complete and no light exits the modulator. When the two paths are in phase constructive interference occurs and the light exiting is equal in intensity to the incident light (minus small system losses). FIG. 1 plots the percentage of transmitted light at different control voltages. The transfer characteristic of the M-Z modulator is a raised cosine. The difference in drive voltage between the minimum and maximum transmission points is called Vxcfx80 since it results in an optical phase change of xcfx80 radians. Vxcfx80 is generally stable, however, the absolute voltage of the optical minimum and maximum drift with time, temperature, and other phenomenon.\nTypically, commercial M-Z modulators are driven with two electrical signals. The first is a data drive signal that modulates the phase difference between the two paths. The second is a bias voltage that controls the average phase difference between the paths. These electrical drive components can be combined externally to the M-Z, and applied to a single electrode set, or can be effectively combined within the M-Z modulator (on separate data and bias electrode sets). This separation of data and bias signals generally simplifies the overall electrical drive system. The bias voltage circuitry corrects for the relatively slow drift due to time, temperature, etc., freeing the high-performance data drive circuitry from addressing these effects. In some M-Z modulators, the data drive signal consists of a plurality of wires that may be driven differentially, redundantly or complementarily. The methods described in this patent apply equally to M-Z modulators with one or more data drive wires.\nSome systems modulate multiple bands of electrical information onto a single optical carrier frequency. These systems are most efficient when the transmitter has a linear transfer characteristic. This is best achieved with the M-Z modulator by setting the bias point at a 50% transmission level (labeled xe2x80x98Axe2x80x99 on FIG. 1) and using the data drive signal to vary the transmission symmetrically about this point. Additional linearity can be achieved by limiting the data drive signal variance to less than Vxcfx80.\nOther systems modulate the optical carrier with a single channel of digital electrical information. Increased data transmission rates are accomplished by increasing the electrical clock rate as well as the use of multiple modulators operating on multiple optical carriers of different frequencies (i.e., Wavelength Division Multiplexing). In these systems, linearity of the transmitter transfer characteristic is unimportant. Rather, the ability of the receiver to distinguish optical xe2x80x981xe2x80x99 from optical xe2x80x980xe2x80x99 must be maximized. We will call the ratio of the amplitude of the optical xe2x80x981xe2x80x99 to the amplitude of the optical xe2x80x980xe2x80x99 the extinction ratio. The higher the extinction ratio of the transmitter, the lower the required sensitivity of the receiver. The maximum extinction ratio is achieved by providing a data drive signal variance of Vxcfx80 combined with appropriate biasing. The term xe2x80x981xe2x80x99 level or xe2x80x980xe2x80x99 level is used to indicate the optimal optical amplitude at the transmitter when sending a xe2x80x981xe2x80x99 bit or a xe2x80x980xe2x80x99 bit to the receiver.\nIn prior-art systems, the data drive signal varies symmetrically and the bias voltage is maintained at the 50% level. Many bias voltage regulation schemes are known, one example is U.S. Pat. No. 6,046,838 titled, AUTOMATIC BIAS CONTROL FOR ELECTRO-OPTIC MODULATORS. In general, two prior-art bias methods are known. Both techniques rely on injecting a low frequency control tone into the M-Z modulator, extracting a resulting component of the control tone from the output of the M-Z modulator and adjusting the bias voltage by processing the extracted component and optionally, the original control tone.\nIn one technique, a control tone is modulated onto the bias voltage. The bias voltage is adjusted to maximize the detected control tone component. FIG. 6 shows an example of this system. Laser 700 generates an optical carrier transmitted through M-Z modulator 710. M-Z modulator 710 modulates the carrier by combining data drive signal 760 and bias voltage 770. Bias voltage 770 is generated by combining the output of control tone oscillator 750 and feedback control 730. The output of the M-Z modulator 720 is monitored by photo-detector 720. Feedback control 730 extracts from photo-detector 720 the control tone component that is combined with the output of control tone oscillator 750 to generate the necessary bias adjustment.\nIn another technique, the data drive signal is modulated by a control tone. The extracted control tone is multiplied by the original control tone to generate an error signal. Feedback correction of the bias voltage is done using well-known control techniques like proportional and integrative.\nIncreasing the transmitted data rate faces two difficulties. First, the shorter electrical wavelength of the modulating signal shortens the practical active length of the M-Z waveguide causing Vxcfx80 to increase. Second, it becomes increasingly difficult to achieve a particular voltage swing at higher electrical clock rates. These two difficulties combine to lower the actual data drive voltage swing below Vxcfx80, reducing the extinction ratio. When the data drive voltage swing is less than Vxcfx80, the system is said to be underdriven. One response to this effect is to increase the sensitivity required of the receiver. However, there are limits to receiver sensitivity due to a number of factors, including electrical noise in the photo-diode, optical noise injected by optical amplifiers in the transmission path and others. Consequently, it is desirable to maximize the extinction ratio of an M-Z modulator when the data drive voltage swing is less than Vxcfx80. The optimal extinction ratio for an underdriven M-Z modulator is achieved when the bias point is not at the 50% transmission level of the prior-art.\nA goal of the invention is maximize the extinction ratio when the M-Z modulator is underdriven.\nThe amplitude of control tone modulation of the zero and one levels of the M-Z control voltage is independently controlled. Detecting the results of the asymmetric control tone modulation allows the system to better control the extinction ratio.\nIn a preferred embodiment, all of the control tone modulation is placed onto the zero level. The system adjusts the M-Z bias point to maximize the extinction ratio."}
{"text": "Not Applicable\n1. Field of the Invention\nThis invention pertains, generally, to very lightweight safety and protective garments for working animals. The garments provide comfortable protection to such animals against insect (primarily biting flies) harassment as well as possessing high visibility, illuminative and identification indicia on any individual garment. Specifically, the invention provides an ensemble, consisting in an unitary or assemblage of garments that are constructed so as to cover, or drape, discrete portions of a working animal\"\"s body and which (independently) provide protection against biting insects, while presenting, for use during times of low or obscured visibility, superficial illuminative indicia for identification.\n2. Discussion of Relevant Art\nSeveral documents have been found that disclose various attempts to provide animal clothing, or garments, that would afford the general characteristics stated above, namely safety and fly-resistive protection; however none appear to show high visibility, illuminating apparatus. The most relevant of these documents have, as they appear to the instant inventor, both advantages and limitations as hereinafter discussed.\nIn U.S. Pat. No. Des. 374,315, issued Oct. 1, 1996 for PROTECTIVE GARMENT FOR CANINES, there is shown a complete cloaking for a dog that covers the animal from crown to tail, stopping just short of covering the tops of the paws. The garment appears to be for warmth, but such a determination of utility is beyond the scope of the design patent. There is no showing of illuminative indications or paraphernalia.\nThe garment morphology shown in U.S. Pat. No. Des. 372,563, issued Aug. 6, 1996 for PROTECTIVE DOG GARMENT, is similar to one of the varied canine coverings of the instant invention. Again, with the design patent, actual utility is unspoken and there exists some ambiguity whether this is a garment for a xe2x80x9cprotectivexe2x80x9d (i.e., protection) dog or a protective garment for a dog. The instant invention\"\"s intent is not served by the xe2x80x9cleggingsxe2x80x9d shown in this patent and in U.S. Pat. No. Des. 374,315; and there appears in neither the above-mentioned paraphernalia.\nA REVERSIBLE HEAT-REFLECTIVE PET GARMENT is disclosed in U.S. Pat. No. 6,089,194, issued Jul. 18, 2000. This article covers the chest, flanks and back of a dog with a reversible cover that presents, outwardly, an alternate heat-reflective or water-repellant surface. The garment resembles yet another of the garment variations entertained by the present invention, but because of its non-porous composition, it is limited to use in only extreme hot/cold or wet climes.\nSimilar to the article of U.S. Pat. No. 6,089,194 is that disclosed in U.S. Patent No. 5,060,458, issued Oct. 29, 1991 and entitled: PROTECTIVE DOG COAT. This is a relatively heavy coat, for a dog of any type, which affords warmth and the addition of highly visible coloring. There is shown no other means for either identification or comfort.\nOne of the most relevant article found, relative to the instant invention, is seen in U.S. Pat. No. 6,128,891 (\"\"891), entitled: PROTECTIVE HORSE MASK. Although confined to but a portion of a horse\"\"s head, this device, by use of a lightweight mesh fabric and fleece, provides face, ear and nostril protection against both sun and noisome insects. Although the device is perceived to be well suited, to its task, the patent falls short of suggesting that the stated protection be applied to the preponderance of the animal\"\"s body. The mask also lacks means for rendering the animal highly visible.\nAlthough coverings for animals are quite plentiful, the instant inventor has not been apprised of a garment for working animals, particularly horses under saddling gear, or dogs, that possesses lightweight, extensive covering, which affords to the animal excellent ventilation, optional fly protection and especially high visibility during low-light conditions.\nSeveral other disclosures show forms of protective garments for equines and canines that will be listed in the References Cited. Most show breathable fabrics or turnout garments and particular shapes; but all fall short of the salient aspects of the instant invention, which can be realized in one or a plurality of safety garments.\nThe U.S. Pat. No. 6,128,891, issued to McMahon on Oct. 10, 2000 for a PROTECTIVE HORSE MASK is incorporated by reference for its showing of a partial head covering for a horse, said covering consisting in the combination of lightweight, fine mesh and fleece fabrics\nMost terms used herein are to be taken as having their customary English meaning. When different or secondary meanings may be applied, the inventor has, with their first usage, employed quotation marks and given their intended definition in parentheses. A few terms, however, are to be read with the following meanings (especially in the claims):\nbreathablexe2x80x94having capability of ventilation or being evapotranspirative;\nensemblexe2x80x94an assemblage of parts (as in a set of clothing);\nevapotranspirativexe2x80x94having the properties of liquid absorption and evaporation;\nfluorescencexe2x80x94luminescence that may persist after removal of excitation;\nilluminati(on)(ve)xe2x80x94practically speaking, any kind of light or light producing (def. illuminative).\niridescencexe2x80x94brilliance or high reflectance of light;\nluminescencexe2x80x94emission of visible light after non-thermal energy stimulation; and\nsaddling gearxe2x80x94the saddle, including cinching and accouterments for attachments.\nThe instant inventor has overcome deficiencies of the early art by providing an ensemble of illuminative/illuminated, well ventilated, insect-resistive garments for working animals, such as horses (especially under saddling gear) or field dogs. Dual safety objectives are satisfied by this invention: first, prevention of animal harassment by biting insects such as deer/black/horse flies, which could cause a horse to stumble or a dog to turn into line of fire; and second, identification of the animal (and rider or master) during conditions of low ambient light. Hereinafter, the descriptions of the invention and its maximum benefits will be described, principally, in respect of the horse, donned with saddling gear (or under saddle): however, the invention is adaptable for use by most equine and canine creatures, while at task or leisure (i.e., at pasture or roaming).\nThe first of these objectives is attained, in varying degree, by cloaking portions of the animal in a single garment or a plurality of garments, the latter ensemble being preferred, but not required. Garments of the ensemble are constructed of lightweight, breathable mesh fabrics of the type that are used for their ventilation or moisture-wicking properties. Moreover, the fabric used by the inventor is produced, and available off the shelf, in iridescent colors. Under fair weather and light (normal) ambient conditions, the ensemble provides: a hood (head covering); a hood-connected cape (neck covering), which descends to the withers and is attached to the saddling gear; and a body drape, consisting in a single garment attached, to saddling gear, so as to cover the horse\"\"s flanks, back and rump laterally and rearward of the saddling gear. The garment fabric is of a mesh small enough to frustrate the common biting flies, yet large enough to xe2x80x9cbreathxe2x80x9dxe2x80x94in the preferred embodiment, having foramen size from about 0.5 mm to about 1.5 mm. The bonnet portion (head cover) of the ensemble stands off the eyes of the animal, thereby providing to them an insect-and light debris-resistive guard or barrier protection. Under such (normal) conditions, both safety objectives of the invention are met when the iridescent fabric is used.\nTo achieve the second objective, under less than optimum visibility conditions, there is provided, on discrete sites of the garments, a series of illuminative strips or patches. Several fabrics of luminescent, fluorescent or otherwise highly reflective character are readily available, on the market. These indicators, or illuminative indicia, are affixed by sewing/stitching, snap fasteners or, like the preferred method of connecting the above garments, by use of materials bearing hooks and loops, as in the product Velcro(copyright)."}
{"text": "As can be seen by reference to the following U.S. Pat. Nos. 2,734,567; 2,805,632; 2,948,083; and 4,841,686; the prior an is replete with myriad and diverse gutter screen assemblies.\nWhile all of the aforementioned prior art constructions are more than adequate for the basic purpose and function for which they have been specifically designed, these patented structures are deficient either from the stand point of either the simplicity or complexity of their respective constructions.\nIn the instance of the simpler constructions, they are normally characterized by a flimsiness and fragility that does not provide sufficient rigidity to the screening element to guarantee a long useful life; whereas, in the instance of the more complex structures, they are much more difficult to manufacture and therefore more expensive bringing into effect the cost versus benefit ratio for the consumer.\nAs a consequence of the foregoing situation, there has existed a longstanding need among consumers for a simple, rugged, and inexpensive gutter guard apparatus that incorporates structural innovations and improvements that are not found in the prior art constructions, and the provision of such a construction is a stated objective of the present invention."}
{"text": "There is a great variety of tools for picking up and manipulating small objects, arranging for example from the simplest tongs and magnetized screwdrivers to extremely complex remote control arms. Many of these tools, such as tongs or pliers, require that the operator maintain his grip on the work piece being handled. Others, such as certain forceps and locking pliers, can be locked in a closed position on the piece being handled. Many of the locking manipulative tools are relatively large, clumsy, otherwise not particularly well-adapted for manipulation of small electrical components, or too expensive for the task. As background, see for example U.S. Pat. No. 4,605,256 which is directed to a locking tool for manipulation of electronic components, and U.S. Pat. No. 3,896,533, which is directed to a tool for inserting and removing circuit components.\nElectrical outlets are frequently crowded environments in which it can be dangerous or difficult to handle components with the fingers. Accordingly, it is desirable to have tools which are small or portable, easily operable and manipulated by one hand, and which can grasp electrical components firmly, capturing them so that they can be manipulated in any orientation for installing and testing or removing electrical outlets from junction boxes. The tool of the present invention admirably fulfills these requirements. It is simple, easily operable by one hand, can be interconnected with the work piece being handled, and is able to capture and retain the electrical work piece being handled until the operator chooses to disconnect the tool from the work piece."}
{"text": "The present exemplary embodiment relates to solar charging protocols for electric or hybrid vehicles. However, it is to be appreciated that the present exemplary embodiment is also amenable to other similar applications.\nElectric vehicles are powered by an electric motor to which electricity is provided by a group of batteries. Operation of the motor depletes energy stored in the batteries. Electric vehicles are typically recharged from an external power source. For example, the electric vehicle can be recharged at a home or office location by being plugged into a standard outlet. Also, commercial fast charging stations are becoming more commonly available where a higher current charge can be delivered.\nA hybrid vehicle operates using both hydrocarbon fuel and electric power. A conventional engine is fueled by the hydrocarbon fuel while an electric motor is powered by a battery. The engine may operate a generator which charges the battery at times when the full power of the engine is not needed to propel the vehicle. A plug-in hybrid is a hybrid vehicle in which the driver has the option of plugging the vehicle into an exterior electric power source when it is parked so that the battery does not have to be charged by the engine.\nSolar vehicles, as referred to herein, include electric and hybrid vehicles which have one or more solar panels on the body to provide part of the electricity for the electric motor and/or for charging the batteries and further to any vehicle including one or more solar panels that provide electrical power to a vehicle accessory, such as a radio or the vehicle's heating, ventilating and air conditioning system.\nAn exemplary electric vehicle including solar panels is depicted in FIG. 1. Vehicle 10 includes storage batteries 12 mounted within the vehicle. A plurality of solar panels 14 are located on the hood and roof to convert incident solar radiation into electrical energy. The solar panels 14 are electrically connected to the storage batteries 12 and are operative to supply electrical current thereto for recharging. In general, the panels are preferably as large as practical to help create a large total area of solar panels. In practice, this may mean that parts of the panels extend over air rather than the vehicle body.\nFurthermore, notwithstanding the reference to solar panels, other types of systems devised to convert solar radiation into electricity may be suitable for use in the presently disclosed embodiment. Moreover, the present disclosure is suited to any type of apparatus attached to a vehicle that generates electrical energy from solar radiation. Accordingly, the use of the phrase “solar panel” throughout this disclosure is intended to encompass all such apparatus.\nA typical car belonging to an individual is parked most of the time. Therefore, if significant sun exposure can be provided while parked, solar charging can provide a substantial portion of the required energy. In the case of an electric vehicle, the solar vehicle would likely also be a plug-in, so if sunlight is unavailable for any reason (weather, parked underground etc.) the battery can be charged from grid power. In the case of a hybrid vehicle, the battery of the solar hybrid can be charged by the solar panels and by the engine and perhaps also as a plug-in.\nAs compared to a solely plug in electric car, a key advantage of a solar vehicle is that it can be more easily charged at a location away from the operator's home. By utilizing an integrated vehicle controller and GPS system, and/or a smart phone, it is possible to provide guidance to the vehicle operator to maximize vehicle charging via solar energy."}
{"text": "Electronic compasses are not new in the art. Both flux gate and Hall effect circuitry have been used to create suitable directional displays for a variety of navigational tasks.\nFor example, U.S. Pat. No. 4,402,142, issued to Robert C. Dinsmore on Sep. 6, 1983, describes an electronic compass for automobiles utilizing a series of Hall effect sensors. However, the circuits used in that invention have proven to require significant amounts of current, rendering the devices relatively impractical for portable use. Because of the current demands of the Hall effect sensors, amplifiers, Schmitt triggers and display drivers, present Hall effect devices typically consume current at the rate of approximately 200 milliamperes per hour. Extremely large batteries would be necessary to power these devices for extensive periods of time.\nThe within invention eliminates substantial current consuming portions of the typical Hall effect or flux gate-type compass, adds a \"chopper\" circuit which effectively shuts the compass electronics off during most of its operating cycle, and a vibration sensing device which shuts off electrical power to the compass when it is motionless for a specified time."}
{"text": "Retinoic acid (vitamin A acid, abbreviation: RA) is an essential substance to the growth and life support of humans and other mammals. It has been known that retinoic acid acts as a morphogenesis factor in ontogenesis and functions variously in the differentiation and proliferation of adults. For example, it has been known that the acid participates in the cornification, formation of hairs, functions of sebaceous glands, and so on with respect to the epidermis, in the metabolism of bones and cartilages with respect to the connective tissues, in the regulation of immune functions with respect to the immune system, in the differentiation of nerve cells with respect to the nervous system, in the differentiation and proliferation of blood cells with respect to the hemic system, and in the secretion of thyroid hormones, parathyroid hormones and so on and the regulation of the functions thereof in target organs, thus taking part in the mineral metabolism and the basal metabolism. These various physiological actions of retinoic acid are exhibited by directly controlling gene expression through retinoid receptor (RARs, RXRs) family present in cell nuclei. With respect to retinoic acid, there are not only deficiencies but also excesses thereof such as abnormality in cornification, depilation, metabolic disorder of bones and cartilages, and so on. Further, the abnormality of retinoid receptors has recently been found in acute promyelocytic leukemia, head and neck squamous cell carcinoma, lung cancer and so on, and the participation of retinoic acid in the sideration and evolution thereof has been reported.\nIn order to elucidate detailed mechanisms of these various actions of retinoids and to find the possibility for clinical application thereof, it has great significance to develop compounds antagonistic against retinoids. Although TD-550 and TD-560 (Cell Biol. Rev., 25, 209(1991)) and Ro41-5253 (Proc. Natl. Acad. Sci., U.S.A., 89, 7129 (1992)) have already been known as compounds antagonistic against retinoids, they are thought to be poor in both the ability to bind RARs and antagonism against retinoids.\nMeanwhile, RARs and RXRs are known as retionid receptors, which are members of steroid/thyroid receptor superfamily present in cell nuclei. Known receptors belonging to this superfamily include estrogen receptors (ER), thyroid hormone receptors (TR), vitamin D.sub.3 receptors (D.sub.3 R) and steroid hormone receptors. With respect to RXRs, there are .alpha.-, .beta.- and .gamma.-subtypes, and the ligand thereof has recently been identified with 9-cis RA. Further, it has been found that RXRs have the physiological property of forming heterodimers together with RXRs, TR, D.sub.3 R or other receptors. Thus, it is being elucidated that RXRs act synergistically with their respective inherent ligands to take great part in the expression of the functions of retinoic acid, vitamin D.sub.3 or thyroid hormones through such heterodimers. In order to elucidate detailed mechanisms of these various actions of RXRs and to find the possibility for clinical application thereof, it has great significance to develop compounds binding to RXRs.\nIn view of the above actual circumstances, the inventors of the present invention have intensively studied to find that mono- or polyenic carboxylic acid derivatives which will be described below exhibit agonism for RXRs and are useful as drugs. As the prior art, although JP-A-2-76862 and EP 0568898 disclose monoenic carboxylic acid derivatives and polyenic carboxylic acid derivatives these derivatives are different from the compounds of the present invention in both chemical structure and drug effect.\nFurther, the inventors of the present invention have found that heterocyclic compounds described below exhibit extremely high ability to bind RARs and antagonism against retinoids, thus accomplishing the present invention.\nFor example, JP-A-2-240058 discloses heterocyclic compounds which exhibit such a function of agonist and are improved in the adverse reaction due to retinoid excess. However, these compounds are different from the compounds of the present invention in both chemical structure and drug effect."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor fabrication. More particularly, the invention relates to a salicided MOS device and a one-sided salicided MOS device on a semiconductor substrate and a simultaneous fabrication method thereof.\n2. Description of the Related Art\nComplimentary metal oxide semiconductor, (CMOS), devices, for both logic and memory applications, fabricated on the same semiconductor chip, have been commonly used. For example, an embedded dynamic random access memory (DRAM) device includes a memory array and a logic circuit array formed together in a single integrated circuit (IC) chip. This embedded DRAM can thus access a large amount of data at much higher speeds. The embedded DRAM is thus widely used in logic circuitry, to process large amounts of data, such as in a graphic or an image microprocessor. An accomplished embedded DRAM typically includes a logic circuit, a transfer field effect transistor (transfer FET, MOSFET) array, and a capacitor coupled to the transfer FET, wherein the transfer FET serves as a electrode of the capacitor and a selective switch when the transfer FET is selected by a bit line. The voltage status of the capacitor can therefore be read or changed through the transfer FET. One transfer FET typically includes a gate structure and an interchangeable source/drain region on each side of the gate structure. The capacitor is coupled to the interchangeable source/drain region on one side of the gate structure, typically is the source region.\nIn order to obtain lower resistance and increase device speed, a method of performing a self-aligned silicide (salicide) process to form a salicided layer on all exposed silicon surfaces of the gate and source/drain region is proposed. However, the salicide process usually consumes junction depth, causing a shallow junction (also silicide junction), which may cause a charge leakage of the capacitor. For some leakage-concerned products, such as DRAM and CMOS image sensor, non-silicide junction is required for partial positions, for example, the position(s) of source/drain region of MOS connecting capacitor in DRAM and the position(s) of source/drain region of MOS connecting photodiode region in CMOS image sensor.\nFIGS. 1-3 are sectional views of a portion of a semiconductor substrate, schematically illustrating a conventional fabrication process for forming an embedded DRAM. In FIG. 1, an isolation structure 110, such as STI (shallow trench isolation) or FOX (field oxide isolation), is formed in/on a semiconductor substrate 100 so as to create a DRAM active region 102 and a logic active region 104 on the substrate 100. An oxide layer (not shown) and a polysilicon layer (not shown) are sequentially formed on the substrate 100. Then, the oxide layer and the polysilicon layer are patterned to form a first gate structure 116 in the DRAM active region 102 and a second gate structure 118 in the logic active region 104. The gate structure 116 includes a gate 120a and a gate oxide layer 120b; and the gate structure 118 includes a gate 130a and a gate oxide layer 130b. \nIn FIG. 1, using the gate structures 116, 118 as a mask, lightly doped (LDD) regions 140 are respectively formed in the substrate 100 on each side of the gate structures 116, 118 by implantation. Then, a spacer 150 is formed on each sidewall of the first gate structure 116 and a spacer 155 is formed on each sidewall of the second gate structure 118. Then, interchangeable source/drain regions 160, 162, 164, 166 are respectively formed in the substrate 100 on each side of the gate structures 116, 118 by another implantation.\nIn FIG. 2, a photoresist layer 210 is formed on part of the substrate 100 to cover the interchangeable source/drain region 160 which will couple with a capacitor in subsequent process. Then, a self-aligned suicide (salicide) process is performed to form a silicide layer 220 on the gates 120a, 130a and the interchangeable source/drain regions 162, 164, 166.\nIn FIG. 3, the photoresist layer 210 is removed. Then, a capacitor 310 abutting the interchangeable source/drain region 160, such as a trench-type capacitor, is formed in the substrate 100. Number 320 denotes a dielectric film layer, conformal to an inner trench. Number 330 denotes an electrode coupled with the interchangeable source/drain region 160. The substrate 100 serves as another electrode of the capacitor 310.\nAccording to the conventional method, however, since misalignment of the photoresist layer 210 easily occurs in submicron process, the DRAM cell performance cannot be precisely controlled. For example, when the misaligned photoresist layer 210′ covers part of the gate 120a, the silicide layer 220 cannot thoroughly form on the gate 120a, as shown in FIG. 4A. This affects the conductivity of the gate 120a. In addition, when the misaligned photoresist layer 210″ covers part of the interchangeable source/drain region 160, the silicide layer 220 partially forms on the interchangeable source/drain region 160, as shown in FIG. 4B. This may cause junction leakage (silicide junction), thereby decreasing device reliability.\nU.S. Pat. No. 6,277,683 discloses a process for forming a salicided CMOS device and a non-salicided CMOS device on a semiconductor substrate. The method uses only one silicon nitride layer to provide a component for a composite spacer on the sides of the salicided CMOS device, and to provide a blocking shape for the non-salicided CMOS device during metal silicide formation. Nevertheless, this conventional method does not describe the process for forming a one-sided salicided CMOS device.\nU.S. Pat. No. 6,063,706 discloses a process for forming a salicided device and an ESD protective device having no salicide on a semiconductor substrate. The method uses a pad silicon nitride to cover the ESD protective device, and then two-step salicide process is performed on the substrate to form a silicide layer on the surface of a functional region. Nevertheless, this conventional method does not describe the process for forming a one-sided salicided CMOS device.\nU.S. Pat. No. 6,337,240 discloses a process for forming an embedded DRAM. The method implants ions onto a substrate with different dopant concentration to form different thickness gate oxide layers. An insulation layer serving as a mask is formed to cover the DRAM memory region. Then, a salicide process is performed to form a salicided layer on the surface of a logic region. Nevertheless, this conventional method does not describe the process for forming a one-sided salicided CMOS device."}
{"text": "1. Field of the Invention\nThis invention relates generally to semiconductor manufacturing, and, more particularly, to a method and apparatus for performing field-to-field compensation for errors in semiconductor manufacturing processes.\n2. Description of the Related Art\nThe technology explosion in the manufacturing industry has resulted in many new and innovative manufacturing processes. Today's manufacturing processes, particularly semiconductor manufacturing processes, call for a large number of important steps. These process steps are usually vital, and therefore, require a number of inputs that are generally fine tuned to maintain proper manufacturing control.\nThe manufacture of semiconductor devices requires a number of discrete process steps to create a packaged semiconductor device from raw semiconductor material. The various processes, from the initial growth of the semiconductor material, the slicing of the semiconductor crystal into individual wafers, the fabrication stages (etching, doping, ion implanting, or the like), to the packaging and final testing of the completed device, are so different from one another and specialized that the processes may be performed in different manufacturing locations that contain different control schemes.\nAmong the factors that affect semiconductor device manufacturing are effectively initiating and continuing a manufacturing process without significant human interaction, which can cause delays or errors in the manufacturing process. One of the process steps that is adversely affected by such factors is the photolithography overlay process. Overlay is one of several important steps in the photolithography area of semiconductor manufacturing. Overlay control involves measuring the misalignment between two successive patterned layers on the surface of a semiconductor device. Generally, alignment is important to ensure that the multiple layers of the semiconductor devices are connected and functional. As technology facilitates smaller critical dimensions for semiconductor devices, the need for reduced of misalignment errors increases dramatically.\nGenerally, photolithography engineers use results from the analysis of the overlay errors to make updates to exposure tool settings manually. Some of the problems associated with the current methods include the fact that the exposure tool settings are only updated a few times a month. Furthermore, often the exposure tool updates are performed manually. Similarly, improvements in error prevention and correction in other types of semiconductor manufacturing processes are also needed to improve yields in semiconductor manufacturing processes.\nGenerally, a set of processing steps is performed on a lot of wafers on a semiconductor manufacturing tool called an exposure tool or a stepper. The manufacturing tool communicates with a manufacturing framework or a network of processing modules. The manufacturing tool is generally connected to a machine interface. The machine interface is connected to a machine interface to which the stepper is connected, thereby facilitating communications between the stepper and the manufacturing framework. The machine interface can generally be part of an advanced process control (APC) system. The APC system initiates a control script, which can be a software program that automatically retrieves the data needed to execute a manufacturing process. The input parameters that control the manufacturing process are revised periodically in a manual fashion. As the need for higher precision manufacturing processes are required, improved methods are needed to revise input parameters that control manufacturing processes in a more automated and timely manner. Furthermore, wafer-to-wafer manufacturing variations can cause non-uniform quality of semiconductor devices.\nMany times processed semiconductor wafers can have portions that contain errors, such as overlay errors. Often, a small error on a portion of a processed semiconductor can go unnoticed. This small error may cause a persistent bias in that portion, and subsequent processes may add to the error, causing a substantial error. Many times, these small errors can be caused by abnormalities in a stepper lens, errors in the manner the wafer is positioned in a stepper, and the like. Small residual errors, such as errors caused by persistent biases, can cause yield degradation in the manufacturing of semiconductor wafers.\nThe present invention is directed to overcoming, or at least reducing the effects of, one or more of the problems set forth above."}
{"text": "Over the past century the demand for energy has grown exponentially. With the growing demand for energy, many different energy sources have been explored and developed. One of the primary sources of energy has been, and continues to be, the combustion of hydrocarbons. However, the combustion of hydrocarbons is usually incomplete and releases non-combustibles that contribute to smog as well as other pollutants in varying amounts.\nAs a result of the pollutants created by the combustion of hydrocarbons, the desire for cleaner energy sources has increased in recent years. With the increased interest in cleaner energy sources, fuel cells have become more popular and more sophisticated. Research and development on fuel cells has continued to the point where many speculate that fuel cells will soon compete with gas turbines generating large amounts of electricity for cities, internal combustion engines powering automobiles, and batteries that run a variety of small and large electronics.\nFuel cells conduct an electrochemical energy conversion of hydrogen or other fuel and oxygen into electricity and heat. In some cases, conversion of a hydrocarbon fuel to hydrogen can occur within the fuel cell in a process known as “internal reforming.” Fuel cells are similar to batteries, but they can be “recharged” while providing power.\nFuel cells provide a DC (direct current) voltage that may be used to power motors, lights, or any number of electrical appliances. There are several different types of fuel cells, each using a different chemistry. Fuel cells are usually classified by the type of electrolyte used. The fuel cell types are generally categorized into one of five groups: proton exchange membrane (PEM) fuel cells, alkaline fuel cells (AFC), phosphoric-acid fuel cells (PAFC), solid oxide fuel cells (SOFC), and molten carbonate fuel cells (MCFC).\nMost fuel cells typically include four basic elements: an anode, a cathode, an electrolyte, and a catalyst arranged on each side of the electrolyte. The anode is the negative post of the fuel cell and conducts electrons that are freed from hydrogen molecules such that the electrons can be used in an external circuit. The anode includes channels to disperse the fuel gas as evenly as possible over the surface of the catalyst.\nThe cathode is the positive post of the fuel cell, and typically includes channels etched therein to evenly distribute oxygen (usually air) to the surface of the catalyst. The cathode also conducts the electrons back from the external circuit to the catalyst, where they can recombine with the hydrogen ions and oxygen to form water.\nOne of the difficulties encountered with fuel cells is the regulation of excessive flow and/or pressure from the fuel source feeding the fuel cell. It is quite common to have flow or pressure spikes from the fuel source during fuel cell operation. The flow irregularities may be a result of start up, temperature changes, fluctuations in power demands on the fuel cell, or other phenomena. Flow increases and pressure spikes cause operational instability and flood the fuel cell with fuel that cannot be efficiently used.\nAccordingly, there has been some use of flow regulators in fuel cells to reduce pressure and/or flow spikes. The flow regulators often consist of a set of capillaries. However, capillary flow regulators have a small dynamic range, and typical regulators add significantly to the cost of the fuel delivery system. There has also been some use of a bladder and rubber diaphragm to regulate the flow and/or pressure delivered to the fuel cell, but such systems are not sufficiently robust for long-term use and have limited efficacy across varying fuel types.\nThus, there is a need in the art for means of regulating the flow and/or pressure of fuel delivered to a fuel cell."}
{"text": "FIG. 1 shows a block diagram of an access control system described in Japanese Patent Application No. 101,439/1985 filed by the assignee of the present application on May 15, 1985 for an \"Optical Disk Drive Device\".\nAs illustrated, an optical disk 101 has a plurality of recording tracks on the disk. The tracks comprise series of pits disposed at a high density and arranged in a circular or spiral form. The disk 101 is fitted onto a spindle and is rotated by a disk-drive motor 102. The disk-drive motor 102 is rotated under the control of a disk-motor-drive control system 103.\nAn optical head 104 forms a light spot on the optical disk 101. The light spot is moved in the radial direction of the optical disk 101. The optical head 104 comprises a frame 105, a source of light such as a semiconductor laser 106, a collimating lens 107, a polarization beam splitter 108, a .lambda./4 plate 109, an optical-path changing mirror 110, an objective lens 111 which focuses the light beam from the light source 106 onto the medium surface of the optical disk 101 and forms a light spot 115 on the surface, a tracking actuator 112 which provides fine or microscopic movements of the objective lens 111 in the radial direction of the optical disk 101 for accurate positioning of the light spot on a recording track of the optical disk, and a split-photodetector 113 which has a pair of sensor parts adjacent to each other for detecting the return light reflected from the optical disk 101, and produces a pair of electrical signals corresponding to the amount of light received at the respective sensor parts.\nAn addition/subtraction amplifying circuit 114 determines the sum of the outputs from the split-photodetector 113 to produce a sum signal as an information signal (reproduced data signal), and determines the difference between the outputs of the split photodetector 113 to produce a difference signal as a tracking error signal. The tracking error signal is supplied to a track-traverse counter 118 and a speed-detecting circuit 120.\nThe track-traversing counter 118 receives the output signal from the addition/subtraction amplifying circuit 114 and detects the number of tracks traversed by the optical head 104. An output of this counter 118 is supplied to a target-velocity generation circuit 119.\nThe target-velocity generation circuit 119 receives the output signal of the track-traversing counter 118, and, at the time of access, generates a target-velocity signal for the light spot 115. The target-velocity signal is sent to a head-actuator drive control circuit 117.\nThe head-actuator drive control circuit 117 also receives an output signal from a polarity switching circuit 2. On the basis of these signals, the head actuator drive control circuit 117 controls the drive of a head actuator 116, such as a linear actuator.\nWhen the head actuator 116 is driven through the head-actuator drive control circuit 117 to move the optical head 104 in the radial direction of the optical disk 101.\nA speed detection circuit 120 detects the track-traverse speed (the speed with which light spot 115 traverses the tracks on the optical disk 115). The output of the speed detection circuit 120 is fed to the polarity switching circuit 2. The speed detection circuit 120, together with the polarity-switching circuit 2, forms a velocity-detection circuit 124.\nThe polarity-switching circuit 2 receives an output signal from an access-direction command generation circuit 123. Under the control of the output signal of the access-direction command generation circuit 123, the polarity-switching circuit 2 changes the polarity of the output of the speed detection circuit 120. More specifically, the detected speed (a scalar value) is converted into a detected velocity (a vector value) which also shows the direction.\nFIG. 2 is a transfer-function block diagram of the velocity-control system which represents the diagram of FIG. 1. In this drawing, the results of the subtraction between the output signal V.sub.S * of the velocity detection circuit 124 and the output signal V.sub.r of the target-velocity generation circuit 119 is input to a gain compensation circuit 5. The gain compensation circuit 5 determines the frequency band of the velocity control system.\nThe gain compensation circuit 5, as well as a notch filter 122 and a head-actuator drive circuit 6, are built into the head-actuator drive control circuit 117. The notch filter 122 compensates the mechanical resonance characteristics G.sub.L (S) of a block 6 in the head actuator 116.\nThe head-actuator drive circuit 6 is normally of a current drive type, and also contains a drive current detection circuit.\nA block 7 in the head actuator 116 represents a force constant of the head actuator 116. The block 8 represents transfer characteristics. Its input is an acceleration, and its output is a head velocity V.sub.L (the velocity with which the optical head 105 is moved by the head actuator 116). M designates the mass of the movable part, G.sub.L (S) designates the mechanical resonance characteristics of the head actuator, and S represents Laplacean. K.sub.V represents the sensitivity in the velocity detection of the target-velocity generation circuit 119 and the velocity-detection circuit 124, and .tau. represents the track traverse period (period taken for the light spot 115 to traverse a track).\nFIG. 3 shows an example of gain characteristic of G.sub.L (S) in the block 8 in FIG. 2 which is shown to have a large resonance peak at a certain frequency .omega..sub.L (usually in the order of kHz). FIG. 4 on the other hand illustrates the gain characteristics .vertline.G.sub.N (S).vertline. of the notch filter 122 in FIG. 2. G.sub.n (S) is selected so that: EQU .vertline.G.sub.N (S).vertline..perspectiveto..vertline.1/G.sub.L (S).vertline.,\nwhen .omega..sub.N =.omega..sub.L.\nFIGS. 5 and 6 show open-loop characteristics of the system of FIG. 2. FIG. 5 shows a case where G.sub.N (S)=1/G.sub.L (S), while FIG. 6 shows a case where G.sub.N (S).noteq.1/G.sub.L (S).\nThe system operates in the following manner:\nFirst, the disk drive motor 102 is energized through the disk-drive-motor control circuit 103, and the optical disk 101 shown in FIG. 1 begins to rotate. When the rotation speed reaches a predetermined steady value, the tracking actuator 112 is controlled on the basis of a tracking error signal obtained by the photodetector 113 and addition/subtraction amplifying circuit 114. As a result, the light spot 115 begins to follow the center of a track on the optical disk 101.\nAt the time of track access, the number of tracks on the optical disk 101, which have been traversed by the light spot 115, is counted by the track-traversing counter 118, and, at the same time, in accordance with the number of tracks to be traversed to reach the target track, the target velocity, which is output from the target-velocity generation circuit 119, and the track-traverse velocity of the light spot 115, which is detected by the velocity-detection circuit 124, are input to the head-actuator drive control circuit 117, which performs velocity control in which the velocity is reduced to zero as the light spot 115 approaches the the target track.\nOperation of the velocity control system of FIG. 2 at the time of track access will now be described in detail. A velocity deviation signal Ve, which is the difference between the target-velocity signal V.sub.r from the target-velocity generation circuit 119 and the detected-velocity signal V.sub.s * from the velocity detector 124, is transmitted to the actuator drive circuit 6 through the gain-compensation circuit 5 and the notch filter 122. As a result, a certain drive current is applied to the head actuator 116. Due to this drive current, the head actuator 116 begins to move the optical head 104, causing the light spot 115 to traverse the tracks on the optical disk 101.\nIf the track fluctuation velocity, due for example to eccentricity of the optical disk 101, is denoted by Vd, the difference between the velocity of the head actuator 116 and the track fluctuation velocity V.sub.o is detected by the velocity detection circuit 124 as a detected-velocity signal V.sub.S *. This detected-velocity signal V.sub.S * is fed back in a velocity control system and the control is so made that the detected-velocity signal V.sub.S * coincides with the target-velocity signal V.sub.r.\nA loop transfer function (open-loop characteristics) of this velocity-control system from the velocity-deviation signal Ve to the detected-velocity signal v.sub.S * can be expressed as follows: ##EQU1##\nIf it is so designed that the condition G.sub.N (S)=1/G.sub.L (S) is satisfied, the head actuator does not have the resonance frequency in the high-frequency zone, as shown in the upper part of FIG. 5. But if G.sub.N (S).noteq.1/G.sub.L (S) because of manufacturing fluctuations between individual devices, the head actuator may have a resonance frequency in the high-frequency band as shown in FIG. 6. When the peak of this resonance exceeds 0.sub.db, the velocity-control system loses its stability.\nMoreover, because of a certain dead time of the zero-order hold characteristics of the velocity-detection circuit 124, a long delay in phase is observed in the vicinity of the track-traversing frequency.\nBecause of the configuration described above, the conventional optical disk drive device had the following problems:\n(1) The track traverse velocity is detected based on the track traverse period. When the traverse is made at a low velocity, a time delay (dead time) is lengthened, and the velocity control system loses its stability, and because the cut-off frequency of the velocity-control system must be designed low, the velocity deviation will be increased.\n(2) When the light spot 115 traverses drop-out or data address recording portions, such traverse may erroneously be recognized as traverse of tracks, and, in spite of the slow movement of the light spot, the velocity detection circuit 124 erroneously operates as if the track traverse velocity were high. The result is that the velocity control system is disturbed.\n(3) The head actuator 116 typically has a high mechanical resonance at the frequency of several kHz. In order to eliminate this phenomenon, a notch filter 122 is built in the head actuator drive-control system. If, however, there are a plurality of resonance frequencies, a plurality of notch filters need to be provided, and the size of the circuit of the system is therefore enlarged. In addition, where there are differences in the resonance frequency from one device to another, and the resonance frequencies differ from the frequencies of the notch filters, and the velocity-control system is not stable.\n(4) Because the known system does not detect the direction in which light spot 115 traverses the tracks, but simply assumes that the light spot 115 is traversing the tracks in the direction (toward the inner or outer periphery of the disk) in which the access is to be made and determines the required velocity by changing the polarity of the speed. When the speed is low, and the direction in which the track is traversed is reversed due to track fluctuation or disturbances, a positive feed-back is applied to the system, and the optical disk may behave erratically.\nAnother problem of the above-described optical disk drive device is described below:\nAs the above-described optical disk drive device obtains head position control information from the optical disk, the absolute position of the head cannot be reliably detected when operation of the servo-system is disturbed, e.g., in case of application of a large external impact force. In that case, as well as in the case of abnormal operation of the servo-system the head may run out of the proper range toward the center or periphery of the disk and collide with a stopper located at the end of the range of mobility. As a result of this collision, the head receives a blow and can be broken.\nThe same problem relates to magnetic disk devices, or any other equivalent information storage devices.\nAnother example of prior art is explained below with reference to FIGS. 7 and 8.\nFIG. 7 illustrates a block diagram which shows a control system of a known optical disk drive device published in Papers from the General Meeting of the Institute of Electronics and Communications Engineers (IECE) of Japan, 1985, Vol. 7, pp. 7-76 [1170, \"Track Access in a Two-Stage Servo-System\", by Hiroshi Inada and Shigeru Shimono]. FIG. 8 illustrates waveforms of control signals used in connection with the device shown in the block diagram. In these drawings, reference numeral 201 designates an optical disk for recording information, or with information already recorded on tracks which are arranged in the form of equally-spaced concentric circles or in the form of a spiral. Reference numeral 202 designates a light beam by means of which information is transferred to and from the optical disk. A head actuator, e.g., a linear actuator 205 drives a carriage 204 of an optical head 203 and moves the carriage 204 with respect to the optical disk 201 and across the tracks. A tracking actuator 206 is installed on the carriage 204 and carries a focusing lens for the formation of a spot of light beam 202 on the tracks of the optical disk 201. The tracking actuator 206 is moved in the same direction as the linear actuator 205 and can cover only a relatively small, predetermined number of tracks. A split photodetector 207 which detects the information signal transmitted by the optical beam 202 and converts it into an electrical signal and outputs the electrical signal. A sensor of this detector consists of two parts. Each such part of the sensor produces on its output an electric signal corresponding to the quantity of light of the light beam 202 which is incident on this part.\nA subtraction amplifier 211 receives a signal from each sensor part of the split photodetector 207, performs subtraction, and thus detects deviation of light spot of beam 202 from the center of the track on optical disk 201. A velocity detection circuit 212 detects, on the basis of an output signal from the subtraction amplifier 211, the track traverse velocity (the velocity with which the light spot of beam 202 traverses the tracks of optical disk 201 in its movement across the disk). A pulse generation circuit 213 receives signals from the subtraction amplifier 211 and generates a pulse each time the light spot of the beam 202 crosses a track on the disk. A track counter 214 receives a signal corresponding to the track access number N (the number of the tracks that must be traversed to reach the target track from the initial (currently-positioned) track) supplied from outside. The track counter 214 receives pulse signals from the pulse generation circuit 213 and counts down by \"1\" each time a pulse is applied to it, and its count value is the remaining tracks to be traversed to reach the target track. A reference velocity generation circuit 215 receives from the track counter 214 a signal corresponding to the remaining number of tracks, initially determines the reference velocity pattern corresponding to the number of the remaining track, memorizes this pattern, and then sequentially produces on its output the reference velocity signals corresponding to gradual decrease in the number of remaining tracks counted by the counter 214. A velocity error detector 216 receives a reference velocity signal from the reference velocity generation circuit 215 and a light spot velocity signal from the velocity detection circuit 212, and which detects the difference in velocities. An amplifying circuit 217 amplifies an output signal of the velocity error detection circuit 216 and controls the linear actuator 205. A position control command circuit 218 receives signals from the operational amplifier 211, the velocity control circuit 212, and the track counter 214. When on a predetermined track the velocity of the light spot of beam 202 drops below a predetermined value, the position control command circuit 218 produces a position control command on its output. A tracking servo-circuit 219 receives a position control command from the position control command circuit 218 and thus controls operation of the tracking actuator 206.\nThe above-described conventional optical disk drive device operates as follows: Track-access control is comprised of a velocity control mode and a position control mode. In the velocity control mode, the carriage 204 is driven by the linear actuator 205, to cause movement of the light spot in the direction of traverse of the tracks of the optical disk 201. In the position control mode, after the velocity of the light spot of the light beam 202 has been reduced below a predetermined velocity at the predetermined track, the tracking actuator 206 is controlled and the light spot is stopped at the position where the spot coincides with the center of a track on the disk 201 (FIGS. 8A-8C). First, in the velocity control mode, a signal which corresponds to track access number supplied from outside (number N in FIG. 7) is sent to the track counter 214. Because at the very beginning there are no pulses from the pulse generation circuit 213, so the number of the remaining tracks is left unchanged and the generated signal corresponds to this particular number N. Receiving this signal, the reference velocity generation circuit 215 initially determines the reference velocity pattern (FIG. 8A), and then sequentially outputs reference velocity signals in accordance with the number of remaining tracks as counted by the track counter 214. The reference velocity signal and the light spot velocity signal, which is produced by the velocity detection circuit 212, are input to the velocity error detection circuit 216 where both signals (i.e., of detected and reference velocities) are compared. The difference is amplified by the amplifier 217, and the amplified signal is used to control the velocity of the linear actuator 205. In accordance with the reference velocity pattern, the linear actuator 205 makes acceleration up to a predetermined number of tracks, the velocity is then stabilized until a predetermined number of tracks is reached, when the deceleration is made.\nIn this way, the light spot of beam 202 moves across the optical disk to reach the target track. When the light spot of beam 202 traverses a track, the quantity of light reflected from the optical disk 1 will change. As the sensor of photodetector 207 consists of two parts, the quantity of light reflected onto each sensor part also will vary. The light reflected onto the sensor is converted to electrical signals which correspond to the amount of light received by the sensor and which are output from the sensor parts. The output signals from the sensor parts are input to the subtraction amplifier 211 performs subtraction to produce a difference signal as shown in FIG. 8D. In this difference signal waveform, a zero point of each cycle corresponds to the moment when the center of the track on optical disk 1 coincides with the center of the light spot of beam 202. The velocity detection circuit 212 receives the output difference signals from the subtraction amplifier 211 and detects on the basis of these signals the track traverse velocity. The pulse generation circuit 213 generates pulses, for example, at the moment of each cycle when the difference signal waveform of the output from the subtraction amplifier 211 passes through zero. Each such pulse is used as a signal indicating that the light spot of beam 202 crossed the track. The pulses are supplied to the track counter 214. The position control command circuit 218 receives output signals of the subtraction amplifier 211, the velocity detection circuit 212, and the track counter 214. If at the moment of arrival of the light spot at a position with a predetermined number of tracks to the target track, e.g., one track to the target one, the velocity is below a predetermined value, the position control command circuit 218 will issue an output command which will switch the system to the position control mode.\nIn the position control mode, the tracking servocircuit 219 receives the output signals of the position control command circuit 218 and the subtraction amplifier 211, and controls the tracking actuator 206 referring to the phase of the difference signal waveform from the subtraction amplifier 211. When the center of the target track of the optical disk 201 coincides with the center of the light spot of beam 202, the tracking actuator stops. Thus, pull-in into the track is completed. The light spot of beam 202 follows the target track, and recording and reproduction of information is conducted.\nIn the known optical disk drive device of the type described above, the carriage 204 is driven by the linear actuator 205, and when the target track is reached and operation of the linear actuator 205 and tracking actuator 206 is switched from the velocity control mode to the tracking control mode, the detected speed may be disturbed either by defects in optical disk 201, or by sudden deviation in the actual velocity due for example to external forces. As a result, the pull-in by the tracking servocircuit 219 may not be achieved. In such a case, the system may behave erratically, unless an external position or velocity scale is provided."}
{"text": "The invention relates to a security or tamper-indicating tag, i.e. a tag that indicates that it has been removed after an initial application to an item. More particularly, it relates to a tamper-indicating tag having two panels that adhere to each other, and a strap that encircles a portion of the object being labelled.\nTags having two indicia-receiving panels that adhere to each other are well known. One particular type of two-panel tag has a strap made of the same material as the two tag panels. This strap extends between the tag panels to automatically form a loop about an item, such as a ring, when the two panels are adhered to each other. A similar tag (commonly known as a \"rat-tail\") has two panels foldably connected to each other and a tail of the same material projecting from one end on one of the panels. The opposite, initially free, end of the tail forms a loop when it is placed between the panels before they are adhered to each other.\nSingle-panel tamper-indicating tags are also known. The indicia-receiving panels of such tags have more than one layer. The layers are chosen so that, after the panel is adhered to an item, the layers will delaminate when the tag is removed. This delamination either destroys the tag layer carrying the indicia with which the tag labelled the item, or visibly defaces the indicia layer.\nSuch tamper-indicating tags have many uses. Their most common use involves the permanent placement of a manufacturer's serial number, warranty, or other identification upon an item. Anyone wishing to transfer the tag to another item would find this task impeded by the delamination of the tamper-indicating tag, with some layers being removed while others remain behind, still adhered to the original item.\nThe use of particular interest here is as a price tag, to prevent customers from switching a tag marked with a lower price indicia from one item to another item of higher price.\nBecause such single-panel tamper-indicating tags are adhered to the item, they are not as suitable for use with small or irregularly shaped items, such as rings, for example, as are two-panel tags.\nA further disadvantage exists when a single panel tamper-indicating tag is used for purposes like price marking which ordinarily require complete removal of said tag. The removal of a one-panel tamper-indicating tag necessarily leaves a portion of the tag, including some adhesive, on the item when the tag delaminates upon removal. This residue is undesirably messy and difficult to remove."}
{"text": "The present invention relates to the field of antennas, and more particularly, to the field repair and replacement of phased array antennas.\nFor phased array antennas, such as electronically scanned array (ESA) antennas, there is an emerging requirement to utilize modular arrays, in which standardized units or portions of the antenna (e.g., sub-arrays or a radio frequency (RF) feed network) are replaceable in the field as part of mission support. Driving this requirement is the desire to simplify and reduce the cost of repair or replacement of part of the antenna, for example, by reducing the size and cost of spares. Further, after replacement, the phase and amplitude of the antenna elements of a newly replaced sub-array, or those corresponding to a newly replaced feed network, must be calibrated (a process typically called phase-up). Thus, there is a desire in the art to eliminate the need to remove the entire antenna from the platform and either utilize special test equipment (STE) in the field or return it to the factory for recalibration or phase-up.\nOne conventional approach utilizes near field techniques through the use of a portable RF absorber aperture cover with an embedded horn feeding a network analyzer. The cover is placed over the aperture and a coarse measurement of the phase and gain of the replaced elements is made and used to align the new elements to the rest of the array. Another similar technique has horn antennas mounted on the edges of the aperture and the signals are processed within the system.\nStill another approach is taught in U.S. Pat. No. 5,657,023 issued to Lewis et al., the entire content of which is incorporated herein by reference. Lewis provides for phase-up of array antennas of a regularly spaced lattice orientation, without the use of a nearfield or farfield range. The technique uses mutual coupling and/or reflections to provide a signal from one element to its neighbors. This signal provides a reference to allow for each antenna element to be phased-up with respect to one another.\nReferring to FIG. 1A, as taught in Lewis et al., a line array includes antenna elements 1-5. The sequence begins by transmitting from element 1 as shown in FIG. 1A as transmission T1, and simultaneously receiving a measurement signal R in element 2. A signal T2 is then transmitted from element 3, and a measurement signal is received in element 2. The phase and gain response from element 2 in this case (reception of the transmitted signal from element 3) is compared to that for the previous measurement (reception of the transmitted signal from element 1). This allows the transmit phase/gain differences between elements 1 and 3 to be computed. While still transmitting from element 3, a receive measurement is then made through element 4. The differences in receive phase/gain response for elements 2 and 4 can then be calculated.\nTo finish the example depicted in FIG. 1A, a signal T3 is transmitted from element 5 and a receive signal is measured in element 4. Data from this measurement allows element 5 transmit phase/gain coefficients to be calculated with respect to transmit excitations for elements 1 and 3.\nThe result of this series of measurements is computation of correction coefficients that when applied allow elements 2 and 4 to exhibit the same receive phase/gain response. Further, additional coefficients result that when applied, allow elements 1, 3 and 5 to exhibit the same transmit phase/gain response. Typically, the coefficients can be applied through appropriate adjustment of the array gain and phase shifter commands, setting attenuators and phase shifters.\nIn a line array of arbitrary extent, the measurement sequences of transmitting from every element and making receive measurements from adjacent elements continues to the end of the array. Thus the calibration technique can be applied to arbitrarily sized arrays. Receive measurements using elements other than those adjacent to the transmitting elements may also be used. These additional receive measurements can lead to reduced overall measurement time and increased measurement accuracy.\nFor an odd element receive phase-up the second series of measurements is aimed at phasing up the odd numbered elements in receive and even numbered elements in transmit. These measurement sequences are similar to those described above for the even element phase-up, and are illustrated in FIG. 1B.\nFirst, a transmit signal from element 2 provides excitation for receive measurements from element 1 and then element 3. This allows the relative receive phase/gain responses of elements 1 and 3 to be calculated.\nA transmit signal from element 4 is then used to make receive measurements from element 3 and then element 5. This allows the relative receive phase/gain response of elements 3 and 5 to be calculated. Also, the relative transmit response of element 4 with respect to element 2 can be calculated. All of the coefficients can then be used to provide a receive phase-up of the even elements and a transmit phase-up of the odd elements.\nTo complete the overall phase-up utilizing conventional practices, the interleaved phased-up odd-even elements need to be brought into overall phase/gain alignment. Coefficients are determined, which, when applied, achieve this alignment.\nHowever, in accordance with the technique described in Lewis et al. each individual antenna element is measured and calibrated, which can be time consuming and energy wasting."}
{"text": "There are a number of existing methods and apparatuses which are capable of providing for real-time monitoring of a patient's vital statistics (“vitals”). These apparatuses include electrocardiogram recorders, heart rate monitors, blood pressure monitors, electroencephalograph apparatus, pulse monitors, oximeters, carbon dioxide meters, thermostats, scales, maternal uterine activity monitors, and various other non-invasive medical instruments.\nA major concern with current noninvasive medical instruments is that they are often bulky and require an excessive number of cables in order to report the measurements to a computer or physician. Everyone is familiar with the image of a hospital patient excessively wired to many machines surrounding or attached to the patient's bed. Access to such patients is difficult during normal treatment, but becomes a real problem in emergency situations where rapid evacuation may be required. In emergencies, non-ambulatory or bed-ridden patients must be quickly and safely evacuated from hospitals, a situation in which patients wait (often for hours in a parking lot) for transportation to another hospital. In such emergencies, because of the issues mentioned above, most current medical instrumentation cannot be quickly and effectively packed up to preserve substantially uninterrupted patient monitoring. Similarly, in the case of home health care, it would be highly desirable to simplify and/or create more user-friendly monitoring methods by reducing the amount of equipment and creating a self-contained, maintenance-free monitoring device.\nNumerous improvements in the medical field have been made to reduce the number of monitors physically attached to a patient, and the size and number of devices in a hospital room, by integrating certain sensor devices into the existing bedding of a patient. The following patents and patent publications, which are hereby incorporated by reference in their entirety herein, disclose a number of contactless, non-invasive patient monitoring methods.\nA mattress pad disclosed in U.S. Patent Publication No. 20070149883 to Yesha has at least two pressure-sensitive piezoelectric sensors positioned in a rigid pad beneath the patient's mattress. The mattress pad includes a processor to receive successive sensor measurements and calculate heart and respiration rates, which are determined by subtracting the pressure signals corresponding to the upper body and the lower body of a patient and mathematically determining the maximum difference of signals between each group of sensors. The heart and respiration rates are then transmitted by a cable to existing patient monitoring equipment. This system, however, requires that a rigid pad be properly disposed and positioned beneath the mattress while physically connected to an auxiliary device, restricting both the mobility and versatility of the monitoring system.\nA mattress with integrated piezoelectric sensors disclosed in U.S. Pat. No. 7,652,581 to Gentry has a passive sensor, or sensor array, in the mattress pad that supports continuous monitoring of a patient's physiological condition in a hospital setting. The external processor receives sensor data, either by wired or wireless communication, from mattress pad sensors, and processes the sensed data into a form that is usable by a physician, nurse or other user. As in Yesha, this is not a self-contained monitoring mattress and still requires auxiliary equipment to receive and process sensor data.\nYet another monitoring mattress pad apparatus is disclosed in U.S. Pat. No. 7,164,941 to Misczynski. This document discloses a contactless electromagnetic inductance device which collects cardiac activity signals to evaluate patient sleep. As in Yesha, the system is not self-contained and also does not permit storage of patient data.\nThere are numerous other patents and published patent applications which employ a pad or embedded mattress sensor coupled to an auxiliary device. One apparent disadvantage is the lack of a self-contained system and the reliance on or requirement of an auxiliary device. Another disadvantage is the lack of embedded storage for patient identification and other data, such as a patient's chart information.\nThus, what is needed is a self-contained Smart Mattress which is capable of monitoring vital statistics of a patient in real time, analyzing data using an embedded processor, storing patient identification and medical information, producing an electronic medical report, and communicating pertinent data to a caregiver or computer server using wireless technology, such as Bluetooth® Technology. Due to the self-contained, portable nature of the Smart Mattress and Smart Mattress technology, it is also ideal for use with an emergency evacuation mattress system for hospital patients and other bed-ridden patients."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus for starting and operating a discharge lamp, using an inverter circuit.\n2. Description of Background\nAn apparatus for starting and operating a discharge lamp, equipped with an inverter circuit has generally been used so as to achieve a small-sized unit with an improved efficiency. Japanese Patent Disclosure (KOKAI) No. 57-141895 discloses a conventional apparatus for starting and operating a discharge lamp. According to the invention, an AC power source is connected to a rectifier circuit including a rectifier and a smoothing capacitor, the output of the rectifier circuit being connected to an inverter circuit. The inverter circuit includes a parallel resonance circuit of an oscillation transformer and capacitor connected in parallel with a primary winding of the oscillation transformer. The resonance circuit is driven by a switching transistor. Between the base and the emitter of the transistor, a series circuit of a coupling capacitor and feedback coil magnetically coupled to the oscillation transformer is connected as a feedback circuit, via a diode. A series circuit of a diode and resistor is connected in parallel with the series circuit of the coupling capacitor and feedback coil, and a trigger resistor is connected in series with both the series circuits.\nA discharge lamp is connected to a secondary coil of the oscillation transformer, via a current-limiting coil, and a starting circuit is comprised of a series circuit of diodes and silicon symmetrical switch element, and is connected in parallel with a discharge lamp.\nAccording to the aforementioned apparatus for starting and operating a discharge lamp, since the starting circuit includes a series circuit of diodes and a silicon symmetrical switch element, a high-voltage pulse to be generated is regulated by the diode characteristic. It is, therefore, not possible to obtain an adequate high-voltage pulse."}
{"text": "The present invention relates to a roof assembly for a building. The present invention also relates to a building that includes a roof assembly."}
{"text": "Integrated circuits (ICs) may be designed to include or not include IEEE 1149.1 boundary scan circuitry and interface. The benefits of including 1149.1 in an IC include but are not limited to; (1) boundary scan testing of interconnects between ICs on a substrate, (2) testing of the IC, (3) debugging circuits within the IC and (4) programming circuits within the IC. The drawbacks of including 1149.1 in an IC include; (1) the requirement of a dedicated 4 pin test interface which increases the IC package size and a timing penalty on the ICs inputs and outputs due to the boundary scan cell multiplexers. Most large digital ICs, such as CPUs, DSPs and ASICs, include 1149.1 for boundary scan testing and even more importantly to enable debug, trace and emulation of embedded core circuits within the ICs. However, 1149.1 is not widely used in memory, analog and mixed signal ICs due to the above mentioned increase in package size and decrease in input and output performance. Memories in particular are resistant to using 1149.1 due to JDEC's memory pin out standardization.\nFIG. 1 illustrates an example integrated circuit die 102 that includes a functional circuit 104 that has external functional inputs (FIN) 106 and external functional outputs (FOUT) 108. Circuit 104 may have bidirectional signals as well, but for simplicity only inputs and outputs are discussed in this disclosure. The circuit 104 may be a digital, analog or mixed signal circuit that performs the functional operation of the die. Circuit 104 may also be a memory circuit, such as but not limited to a double data rate random access memory. As can be seen the die of FIG. 1 does not include 1149.1 boundary scan circuitry.\nFIG. 2 illustrates an example integrated circuit die that includes a functional circuit 104 and 1149.1 circuitry consisting of a test access port (TAP) 204 and associated boundary register (BREG) 206. The TAP has an external interface of TDI, TCK, TMS inputs 212 and a TDO output 214 signals. The TAP responds to the TCK and TMS signals to input data from TDI and output data to TDO. If the boundary register is selected for access it will shift data from TDI to TDO. During normal operation of the die, the boundary register couples the FIN signals 106 to the internal inputs 208 of the circuit and the internal outputs 210 of the circuit to the FOUT signals 108. During boundary scan test mode using the well known 1149.1 Extest instruction, the boundary register isolates the FIN signals 106 from the internal inputs 208 and the internal outputs 210 from the FOUT signals 108. In the boundary scan test mode the boundary register can be operated to capture test data from the FIN signals 106 and update test data to the FOUT signals 108. The boundary scan test mode enables the testing of the FIN and FOUT connections between multiple die/ICs on a substrate.\nFIG. 3 illustrates the TAP 204 of die 202 in more detail. The FIN signals 106 and the 1149.1 input signals 212 define the inputs 318 to the die. The FOUT signals 108 and the 1149.1 output signal 214 define the outputs 320 of the die. The 1149.1 TAP includes, at minimum, a TAP state machine (TSM) 304, an instruction register 306, a Bypass Register 308, the Boundary Register 206 and TDO output multiplexers 310 and 312. While not shown, the TAP may also include other data register between TDI and TDO, including an optional Identification Register which contains a 32 bit code identifying the die. During functional operation, the FIN signals 106 and FOUT outputs signals 108 are coupled to the circuit 104 via the boundary register 206 and buses 208 and 210. The TSM 304 operates according to the well known 16 state transition diagram of FIG. 3A in response to the TCK and TMS input signals to; (1) place the TAP is a Test Logic Reset state, (2) place the TAP in a Run Test/Idle state, (3) perform a scan operation to the instruction register from TDI to TDO, (4) to perform a data scan operation to the Bypass Register 308 from TDI to TDO or (4) perform a data scan operation to the Boundary Register 206 from TDI to TDO. The 1149.1 input interface 212 may include an optional TRST input, shown in dotted line, to reset the TSM and other TAP circuits. If the TRST input is not included, a Power Up Reset (POR) circuit 316 may be used to reset the TSM and other TAP circuits.\nDuring instruction scan operations, the TSM outputs control (CTL) signals to the instruction register 306 and multiplexer 312. In response to the CTL signals the instruction register performs capture, shift and update operations. During the shift operation the instruction register shifts data from TDI to TDO via multiplexer 312.\nDuring data scan operations, the TSM outputs CTL signals to the selected data register 308 or 206 and multiplexer 312. The instruction register output (IRO) bus enables the selected data register and controls multiplexer 310 to couple the TDO output of the selected data register to the TDO output of the die via multiplexer 312. In response to the CTL signals the selected data register performs capture, shift and update operations, except for the Bypass Register 308 which does not have update circuitry. During the shift operation the selected data register shifts data from TDI to TDO via multiplexers 310 and 312.\nDuring manufacturing test of die 202, the inputs 318 and outputs 320 are connected to a tester. The tester operates the inputs 318 and outputs 320 to test the circuit 104 within the die 202. The test may include operating only the FIN 106 and FOUT 108 signals, operating the FIN 106, FOUT 108 and 1149.1 input 212 output 214 signals, or operating only the 1149.1 input 212 and output 214 signals. After testing and found to be good, the die is ready for use within in a system.\nFIG. 4 illustrates three die 402-406 connected to a system substrate 408 via their inputs 318 and outputs 320. The system of this and following examples could be any type of electronic system such as a computer system or a cell phone system. Each die contains a functional circuit 104 and a TAP 204 as shown in FIGS. 2 and 3. The functional circuits 104 of each die provide a different functional operation on the substrate.\nFIG. 5 illustrates the circuits 104 of die 402-406 connected on the system substrate 408. The substrate provides a functional bussing path 502 that connects the FIN 106 and FOUT 108 signals of the die together to enable the die to communicate. The substrate has FIN signals 504 to allow it to input signals from an external device, such as a keyboard, and FOUT signals 506 to allow it to output signals to an external device, such as a display.\nFIG. 6 illustrates the TAPs 204 of die 402-406 connected on the system substrate 408. The substrate provides a serial bussing path 602 that connects the TAPs 204 to externally accessible 1149.1 input 212 and 1149.1 output 214 signals. When an 1149.1 controller is connected to the external 1149.1 input and output signals, the TAPs can be serially accessed to perform test or other operations. One of the most important test operations the TAPs perform is the verification that the FIN 106 and FOUT 108 of each die are properly connected together via the substrate bussing path 502. This test operation is performed by loading Extest instructions into each TAP's instruction register then operating the boundary registers of each TAP to test the connectivity between each die's FIN 106 and FOUT 108 signals.\nFIG. 7 illustrates a device 700 comprising three stacked die 702-706 mounted on a silicon interposer 708. Each die includes a circuit 104 and a TAP 204 as described in FIGS. 2 and 3. Again, each die circuit 104 will typically provide a different functional operation. Die 702 and 704 in this example are designed using through silicon vias (TSV) 710. TSVs are connectivity paths formed between the top and bottom surfaces of the die. TSVs allow input 318 and output 320 signals to flow vertically up and down the die stack. In addition to TSV input and output signal connections from the interposer, the die 702-704 are also connected together locally using input 712 and output 714 signal connections. Interposers 708 are used to provide electrical connections between one surface and another surface. The primary purpose of an interposer is to spread connections from fine pitch contact points on one surface to wider pitch contact points on another surface. In this example, the fine pitch contact points on the bottom surface of die 702 are spread to match the wider pitch contacts points of a system substrate device 700 will be mounted on. Once the device 700 is mounted on a substrate, it receives input signals 318 from the substrate and sources output signals 320 to the substrate.\nBefore device 700 is assembled, each die 702-704 and the interposer 708 are tested to insure the device is assembled with known good die and interposer. As mentioned in regard to FIG. 3 the test may performed by a tester operating some are all of the FIN 106 and FOUT 108 signals of buses 318 and 320, operating some are all of the FIN 106, FOUT 108, and the 1149.1 input 212 output 214 signals of buses 318 and 320, or operating only the 1149.1 input 212 and output 214 signals of buses 318 and 320.\nFIG. 8 illustrate a first example arrangement of how the TAPs 204 of die 702-706 in device 700 may be accessed by a tester. In this example, interposer signal bussing path 802 provides the TDI and TCK inputs to all TAPs, interposer signal path 804 provides a TMS1 input to the TAP of die 702, interposer signal path 806 provides a TMS2 input to the TAP of die 704, interposer signal path 808 provides a TMS3 input to the TAP of die 706 and interposer signal path 810 provides the TDO outputs from all the TAPs. Having unique TMS1-3 inputs for each TAP allows enabling one TAP while the other TAPs are disabled. When disabled a TAP's TDO output is tri-stated to avoid contention on signal path 1206 with an enabled TAP's TDO output. This is commonly referred to as the 1149.1 Star mode of accessing TAPs. To access the TAP of die 702 for a test or other operation, the tester inputs TMS1, TDI and TCK signals to the TAP via busses 804 and 802, and receives TDO signals from the TAP via bus 810. To access the TAP of die 704 for a test or other operation, the tester inputs TMS2, TDI and TCK signals to the TAP via busses 806 and 802, and receives TDO signals from the TAP via bus 810. To access the TAP of die 706 for a test or other operation, the tester inputs TMS3, TDI and TCK signals to the TAP via busses 808 and 802, and receives TDO signals from the TAP via bus 810.\nFIG. 9 illustrates three devices 902-906 connected to a system substrate 908 via their inputs 318 and outputs 320. Each device contains a stack of die 910 with TAPs and an interposer 708 as shown in FIGS. 7 and 8.\nFIG. 10 illustrates the substrate providing a functional bussing path 1002 that connects the FIN 106 and FOUT 108 signals of the devices 902-906 together to enable them to communicate. The substrate has FIN signals 1004 to allow it to input signals from an external device, such as a keyboard, and FOUT signals 1006 to allow it to output signals to an external device, such as a display.\nFIG. 11 illustrates the substrate providing an externally accessible 1149.1 signal bussing path 1102 to the TAPs of the devices 902-906. Assuming the devices are assembled as shown in FIG. 8, the 1149.1 bussing path 1102 would include 9 unique TMS signals. In this example, TMS signals 1-3 would be used to individually access one of the three TAPs of device 902, TMS signals 4-6 would be used to individually access one of the three TAPs of device 904 and TMS signals 7-9 would be used to individually access one of the three TAPs of device 906.\nWhen an 1149.1 controller is connected to the externally accessible 1149.1 bussing path 1102, a selected TAP in device 902 is enabled by one of the TMS1-3 inputs, a selected TAP in device 904 is enabled by one of the TMS4-6 inputs and a selected TAP in device 906 is enabled by one of the TMS7-9 inputs so that they can be serially accessed from the external 1149.1 input bus 212 and 1149.1 output bus. After accessing this first group of serially connected device TAPs, the 1149.1 controller can select a second group of serially connected device TAPs for access using a different set of TMS signals, and so on. Accessing separate groups of device TAPs can be used for performing an 1149.1 Extest operation to verify that the FIN 106 and FOUT 108 signals of each device are properly connected together via the substrate bussing path 1002.\nA first problem with the 1149.1 access approach of FIG. 11 is that the 1149.1 bus 1102 requires a large number of TMS signals that must be routed through the substrate and connected to the multiple TAPs of devices 902-906.\nA second problem with the 1149.1 access approach of FIG. 11 is that standard 1149.1 controllers typically only provide a single TMS signal to support 1149.1 access approaches as shown in FIG. 6. 1149.1 access approaches like that shown in FIG. 11 would require modifying standard 1149.1 controllers to include and operate multiple TMS signals.\nA third problem with the 1149.1 access approach of FIG. 11 is that Extest operations are encumbered by having to individually select different groups of device TAPs to access their boundary register to test the FIN 106 and FOUT 108 connections to the functional bus 1002 of the substrate 908.\nA fourth problem is that Extest operations are lengthened due to having to shift test data through boundary register cells of the local inputs 712 and outputs 714 of each die in the device. For example, when the boundary register of the TAP of die 704 of FIG. 8 is being accessed during an Extest operation to test the die's FIN 106 and FOUT 108 connections to a substrate, the boundary register cells on the die's local inputs 712 and outputs 714 also have to be shifted which adds to the test time. Since the local inputs 712 and outputs 714 of each die were tested after the device was assembled, as mentioned in regard to FIG. 7, they do not need to be tested again when the device is mounted on the substrate. Only the device's FIN 106 and FOUT 108 connections to the substrate need to be tested.\nFIG. 12 illustrate a second example arrangement of how the TAPs 204 of die 702-706 in device 700 may be accessed by a tester. In this example, interposer signal bussing path 1202 provides the TCK and TMS inputs to all TAPs, interposer signal path 1204 provides a TDI input to the TAP of die 702, a local signal path 1208 provides the TDO output of the TAP of die 702 to the TDI input of the TAP of die 704, a local signal path 1210 provides the TDO output of the TAP of die 704 to the TDI input of the TAP of die 706 and interposer signal path 1206 provides a TDO output from the TAP of die 706. In this arrangement all the die TAPs of a device are connected in a daisy-chain and can be serially accessed together by the tester.\nFIG. 13 illustrates three devices 1302-1306 connected to a system substrate 1308 via their inputs 318 and outputs 320. Each device contains a stack of die 1310 with TAPs and an interposer 708 as shown in FIG. 12.\nFIG. 14 illustrates the substrate providing a functional bussing path 1402 that connects the FIN 106 and FOUT 108 signals of the devices 1302-1306 together to enable them to communicate. The substrate has FIN signals 1404 to allow it to input signals from an external device, such as a keyboard, and FOUT signals 1406 to allow it to output signals to an external device, such as a display.\nFIG. 15 illustrates the substrate providing an externally accessible 1149.1 signal bussing path 1502 to the TAPs of the devices 1302-1306. Assuming the devices are assembled as shown in FIG. 12, the 1149.1 bussing path 1102 would provide a serial path through all the device TAPs using a single TMS signal. Signal path 1504 in each device is provided to indicate the daisy-chaining of the device TAPs.\nWhen an 1149.1 controller is connected to the externally accessible 1149.1 bussing path 1502, all TAPs in devices 1302-1306 can be serially accessed from the external 1149.1 input 212 to the external 1149.1 output 214.\nA first problem with the 1149.1 access approach of FIG. 15 is that the IEEE 1149.1 standard has rules that support only one TAP in a device to be implemented in a system. For example a device should only have one TAP instruction register 306, one TAP bypass register 308, one optional TAP identification register and one TAP boundary register 206. As seen in FIG. 15, the system devices 1302-1306 each have multiple TAP instruction registers 306, multiple TAP bypass registers 308, multiple TAP boundary registers 206 and multiple optional TAP identification registers. Therefore the system devices 1302-1306 of FIG. 15 are considered to be non-compliant with the IEEE 1149.1 standard.\nA second problem is that Extest operations are lengthened due to having to shift test data through boundary register cells of the local inputs 712 and outputs 714 of the die within each device. For example, when the boundary registers of the TAPs of the device of FIG. 12 are being accessed during an Extest operation to test the device's FIN 106 and FOUT 108 connections to the substrate of FIG. 15, the boundary register cells on the die's local inputs 712 and outputs 714 also have to be shifted which adds to the test time. Since the local inputs 712 and outputs 714 of each die were tested after the device was assembled, as mentioned in regard to FIG. 7, they do not need to be tested again when the device is mounted on the substrate. Only the device's FIN 106 and FOUT 108 connections to the substrate need to be tested.\nThe following disclosure provides a solution to the above mentioned problems of FIGS. 11 and 15. The solution is based on the concept of improving conventional interposers to include IEEE 1149.1 TAP circuitry."}
{"text": "1. Field of the Invention\nThe present invention relates to an information matching system and its matching method for matching information on the state of equipment to be managed by managing equipment (hereinafter, referred to as management information) with actual state of the managed equipment in communication between terminal equipment connected via a network.\n2. Description of the Related Art\nIn a system having a plurality of terminal equipment via a network, the plurality of terminal equipment consisting of managing equipment for managing the state of each terminal equipment and managed equipment to be managed by the managing equipment, an event notice informing the change, if a change of the state occurred in some managed equipment, is issued from the managed equipment to the managing equipment, and the management information in the managing equipment is updated, thereby performing information matching processing. A change of the state means, by way of example, occurrence of a failure, restoration, modification of various setting in the managed equipment, and detachment and attachment of a new hardware package.\nWhen using TCP/IP(Transmission Control Protocol/Internet Protocol) as a communication protocol among terminal equipment over a network, such a TCP (Transmission Control Protocol) having a resending function is used for the communication of information requiring reliability. By the TCP, however, the communication speed is slow for the performance of the resending processing. Therefore, a use of the TCP in all the communication of information would be a great burden and inefficient. Instead, the UDP (User Datagram Protocol) having no resending function and hence enabling high speed communication is used for the above-mentioned event notice informing managing equipment that the state of managed equipment has been changed.\nIf the above-mentioned event notice informing managing equipment about a change of the state of managed equipment is lost on the way in communication, the management information on the managed equipment in the managing equipment is out of accord with the actual state of the managed equipment. When using the UDP as a protocol, the UDP would not resend an event notice, which is incapable of recovering the disagreement between the management information and the actual state.\nIn order to solve this problem, there is a method of periodically reading out information including histories about all the changes of the state of managed equipment to be managed by managing equipment, by use of the TCP having a resending function. This enables the managing equipment to obtain the lost event notice, even if losing given event notice, thereby to match the management information of the managing equipment with the actual state of the managed equipment.\nHowever, a large amount of data communication in periodically reading out the information including histories about all the changes of the state of each managed equipment, results in increasing communication load of terminal equipment and a network."}
{"text": "A system such as described in the patent literature below is considered as a manufacturing system which is constituted such that work machines including a machine tool are disposed in one line and multiple work machines sequentially perform work on one workpiece. In the system, multiple machine tools in which one machine main body section is mounted on one base are lined up, and each machine main body section is attachable to and detachable from each base. Therefore, the system is highly convenient for maintenance, replacement, and the like of the machine main body section."}
{"text": "The invention relates to a spindle motor for memory devices such as hard disk drives.\nConventional spindle motors for hard disk drives include a base plate or a base flange, a shaft, and a bearing system for rotationally supporting a motor's rotor driven by an electromagnetic force. The bearing system may be structured as a rolling bearing system or a hydrodynamic sliding bearing system.\nIn a conventional spindle motor having a stationary shaft, the shaft is firmly connected to the base plate or the base flange by means of a press-fit connection. To accomplish the connection, the base plate or the base flange features a drilled hole of a predetermined diameter and a predetermined length. One end of the shaft is pressed into and secured within the drilled hole. The outside diameter of the shaft is typically slightly larger than the inside diameter of the drilled hole of the base plate. This surplus is selected so as to ensure a minimal adhesive dimension, i.e., a certain adhesive strength, along the entire length of the joint. On the other hand, the surplus is selected to ensure that a maximum adhesive dimension is not exceeded. An excessive surplus would lead to an unacceptable strain on the components.\nThe greater the surplus and joint length, the higher is the press-in force required to join together the base plate and shaft. If the surplus or the joint length is reduced, adhesive strength may become too small. Increasing the surplus or joint length may result in damage to the components during the press-in procedure."}
{"text": "Frequency spectrum is a scarce resource in modern communications systems, and therefore e.g. the Third-Generation Partnership Project (3GPP) specification group is continuously adding new frequency bands.\nMost modern transceiver solution use direct conversion receiver and transmitter architectures to reduce complexity and minimize power consumption. Such transceivers use local oscillators for frequency generation in the receiver part as well as the transmitter part. The receiver and the transmitter need quadrature local oscillator signals. These are normally generated by dividing the local oscillator signal by two. Thus the local oscillator signal should preferably be operating at twice (or a higher even multiple) the desired receive/transmit frequency.\nThe output from the local oscillators has to be a clean low noise signal to get good quality reception/transmission. This is normally done by locking a voltage controlled LC-oscillator to a crystal oscillator using a phase locked loop. The LC-oscillator needs to cover all desired receive and transmit radio frequencies, i.e. a wide range of frequencies has to be supported.\nA typical LC-oscillator consists of an inductor, a variable capacitor, a resistor representing all losses and a sustaining amplifier. The inductor can be implemented as a large on chip metal structure and the capacitor is made variable to be able to generate various output frequencies. The capacitor can consist of switchable capacitors, variable capacitors or a combination thereof.\nIt has proven to be difficult to generate high capacitor variations. When the available range is in-sufficient one has to use multiple VCO:s centered differently and select the correct one via a multiplexer. This generates bulky solutions, and one might need to add a new VCO for each new frequency band to be supported.\nBy adding a frequency divider at the VCO output to divide the VCO output by two the need for frequency variations of the VCO can be reduced. In principle, it would then be sufficient that the VCO covers a frequency range corresponding to one octave, i.e. with a factor two between the upper and lower endpoints of the VCO frequency range.\nHowever, in practice the VCO center frequency will vary across Process, supply Voltage and Temperature (PVT). Thus if the center frequency has drifted downwards or upwards from its nominal value the complete desired frequency range can not necessarily be covered with a relative VCO range of two. Therefore, the variations of the center frequency impose a need for extra margins at each end of the tuning range. Also to avoid respins of an ASIC, some margin should be added for centering error in the design. Altogether this forces the designer to have ˜12% margins at each end of the tuning range, and the required tuning range for continuous tuning range becomes 2.5 instead of the factor 2 mentioned above. This value is still difficult to achieve with one VCO, and therefore multiple VCOs will still have to be used. Typically, it will take three VCOs to cover a tuning range of 2.5.\nIt might be possible to reduce the number of VCOs by doing a more clever VCO design. For instance, the inductor could also be made switchable. However it will always be difficult to reach a relative tuning range of 2.5.\nAmong the main problems with existing solutions are cost and efficiency. Multiple VCOs occupy a lot of space and the multiplexing will consume current. By designing wide tuning range VCOs the number of VCOs required could be reduced. However, when designing wide tuning range VCOs more switching is required, which introduces resonator loss. This means worse noise performance and higher current consumption. Thus there clearly is a need to reduce the tuning range requirements of the VCO."}
{"text": "This invention relates in general to chucks for holding tools and in particular to a new and useful power operated chuck which includes radially movable clamping jaws.\nSuch power chucks, also termed wedge hook chucks, are known in a variety of designs and have proved satisfactory in practice. To equalize the centrifugal forces produced in the clamping jaws at high speeds and reduce an inwardly directed static clamping force or considerably augment an outwardly directed clamping force, it is known to provide balancing elements which are guided in the chuck body as counterweights and operatively connected to the jaws through hinged levers. The principal drawback of these constructions is that the static clamping forces acting on the clamped workpiece before and after a machining operation are not equal to each other and differ sometimes considerably. That is, as the chuck rotates at high speeds, the clamping force transmitted from the clamping piston through the wedge hooks to the jaws becomes almost counterweighted by the produced centrifugal forces and the elasticity of the coupled parts. Then, a part of the clamping force is produced by the balancing elements which are connected to the jaws through reversing levers. The clamping piston coupled to the jaws through the mating surfaces of the wedge hooks is thereby relieved and displaced, even though very little, by the action of a connected servo-mechanism which continues to be effective. In addition, the centrifugal forces expand the chuck body, so that the bore receiving the clamping piston is enlarged, with the result that the no longer backed-up wedge hooks can yield, which causes a further displacement of the clamping piston. However, after a machining operation is terminated, the chuck body contracts again, whereby the wedge hooks are pressed inwardly and the clamping force on the workpiece is augmented since the clamping piston does not retract, it is held fast in its displaced position. In consequence, the clamping force may increase considerably and often leads to a deformation of the workpiece."}
{"text": "The invention relates to a multi-stage transmission.\nA multi-stage transmission is already known from DE 10 2011 108 024 A1 for a motor vehicle having a transmission input shaft for non-rotatable connection of an internal combustion engine, four planetary gear stages which are disposed one behind the other along a main axis of rotation, and a combined gear shift unit. In this case the gear shift unit is designed as a brake and has a rotatable coupling element as well as a friction-locking unit and a form-locking unit connected in parallel with the friction-locking unit, which are in each case provided to connect the rotatable coupling element fixedly to the housing. One of the planetary gear stages of the multi-stage transmission in DE 10 2011 108 024 A1 is associated with the combined gear shift unit and has a transmission element which is permanently non-rotatably connected to the rotatable coupling element of the combined gear shift unit.\nFurthermore, a multi-stage transmission is known from the generic DE 10 2011 112 235 A1 in which additionally the planetary gear stage associated with the combined gear shift unit comprises a further transmission element which is permanently connected non-rotatably to the transmission input shaft.\nThe object of the invention is in particular to reduce any power loss of the multi-stage transmission without affecting the installation space.\nThe invention relates to a multi-stage transmission for a motor vehicle, having a transmission input shaft for non-rotatable connection to an internal combustion engine, having at least four planetary gear stages which are disposed one behind the other along a main axis of rotation and at least one combined gear shift unit which is configured as a brake and which has a rotatable coupling element as well as a friction-locking unit and a form-locking unit connected in parallel with the friction-locking unit, which are in each case provided in order to connect the rotatable coupling element fixedly to the housing, wherein one of the planetary gear stages is associated with the combined gear shift unit and comprises a transmission element which is permanently connected non-rotatably to the rotatable coupling element of the combined gear shift unit.\nFurthermore, it is assumed that the planetary gear stage associated with the combined gear shift unit comprises a further transmission element which is permanently non-rotatably connected to the transmission input shaft. In this way a transmission element of a planetary gear stage configured as an input stage can be particularly advantageously connected frictionally and/or by form-locking non-rotatably to a transmission housing, so that by securing of the transmission element selectively engageable gears with different torque requirements can be shifted in a loss-optimized and convenient manner. Due to the association of the combined gear shift unit with the planetary gear stage, which has a transmission element permanently non-rotatably attached to the transmission input shaft, an installation space can be utilized better, so that the multi-stage transmission can be of particularly compact design. In this way a power loss, in particular a drag effect, of the multi-stage transmission can be reduced in particular without affecting the installation space, in particular without the provision of an additional installation space. “Additionally” should be understood in particular to mean additional by comparison with a multi-stage transmission which lacks a combined gear shift unit associated with a planetary gear stage configured as an input stage. A “gear shift unit configured as a brake” should in particular be understood as a gear shift unit which is disposed functionally between one of the planetary gear stages and the transmission housing and is provided in order, in a closed state, to connect the rotatable coupling element thereof non-rotatably to the transmission housing of the multi-stage transmission. In this context a “form-locking unit” should be understood in particular to be a part of the combined gear shift unit which for fixed connection to the housing has a tooth system and/or claws which interengage by form-locking in order to produce a non-rotatable coupling, such as in particular a claw-type gear shift unit. In this context a “friction-locking unit” should be understood in particular to be a part of the combined gear shift unit which for fixed connection to the housing has friction elements, such as in particular discs, which are connected to one another frictionally in order to produce a non-rotatable coupling, such as in particular a disc-type gear shift unit. The friction-locking unit and the form-locking unit are preferably functionally disposed between the transmission housing and the rotatable coupling element of the combined gear shift unit. “Non-rotatably connected” should in particular be understood to mean a connection in which a power flow averaged over a complete revolution is transmitted with an unchanged torque, an unchanged direction of rotation and/or an unchanged rotational speed. A “combined gear shift unit” should in particular be understood to be a gear shift unit which for production of a non-rotatable connection has a friction-locking unit and a form-locking unit connected in parallel with the friction-locking unit, wherein the non-rotatable connection can be produced by the friction-locking unit alone, by the form-locking unit alone or jointly by the friction-locking unit and the form-locking unit. A “planetary gear stage associated with the combined gear shift unit” should be understood to be a planetary gear stage which has at least one transmission element which is permanently non-rotatably connected to the rotatable coupling element of the combined gear shift unit. A “planetary gear stage designed as an input stage” should be understood in particular to be a planetary gear stage having at least one transmission element which is permanently connected non-rotatably to the transmission input shaft. “Provided” should be understood in particular to mean specially programmed, equipped, configured and/or disposed.\nAccording to the invention it is proposed that the form-locking unit is disposed radially inside the friction-locking unit, so that a circumferential speed of a form-locking tooth system of the form-locking unit can be kept low. The term “radially” and the term “radius” relate in particular to the main axis of rotation, so that the expression “radially” designates a direction which extends perpendicular to the main axis of rotation. A “form-locking unit disposed radially inside the friction-locking unit” should in particular be understood to be a form-locking unit of which the tooth system and/or claws run around a circular path having an average diameter which is smaller than an average diameter of a circular path on the friction surfaces of the friction elements.\nFurthermore, it is proposed and assumed that the transmission element, which is permanently connected non-rotatably to the transmission input shaft, of the planetary gear stage associated with the combined gear shift unit is configured as a sun gear. As a result a driving torque provided by the internal combustion engine can be advantageously introduced into the multi-stage transmission via the planetary gear stage associated with the combined gear shift unit.\nFurthermore, it is proposed that the transmission element, which is permanently connected non-rotatably to the rotatable coupling element of the combined gear shift unit, of the planetary gear stage associated with the combined gear shift unit is configured as a planetary gear support, so that the planetary gear stage configured as an input stage can be advantageously connected.\nIn particular it is advantageous if the planetary gear stage associated with the combined gear shift unit has a transmission element which is configured as a ring gear and can at least be connected non-rotatably to a transmission element of at least one other planetary gear stage, so that the connection of the planetary gear stage associated with the combined gear shift unit can be further improved. The expression “can at least be connected” should be understood in particular to mean that the connection can either be produced by a gear shift unit or is produced permanently.\nFurthermore, it is advantageous if the rotatable coupling element of the combined gear shift unit can at least be connected non-rotatably to a transmission element of at least one of the other planetary gear stages, so that the combined gear shift unit can act on at least one further planetary gear stage.\nFurthermore, it is advantageous if the multi-stage transmission has a housing end wall on the drive side, wherein the planetary gear stage associated with the combined gear shift unit is disposed immediately adjacent to the housing end wall on the input side with respect to the other planetary gear stages. In this way the planetary gear stage associated with the combined gear shift unit can be disposed on the input side, so that the combined gear shift unit can be associated with the planetary gear stage configured as an input stage without affecting the installation space. A “housing end wall on the drive side” should in particular be understood to be a housing end wall which, by comparison with an axially opposing housing end wall on the output side, is disposed closer to the internal combustion engine axially and/or along a flux of force emanating from the internal combustion engine. The expression “disposed immediately adjacent to the housing end wall on the input side with respect to the other planetary gear stages” should in particular be understood to mean that a planetary gear stage is absent axially between the housing end wall on the drive side and the planetary gear stage associated with the combined gear shift unit. The expression “axially” relates in particular to the main axis of rotation, so that the term “axially” designates a direction which extends on the main axis of rotation or parallel thereto.\nIn an advantageous configuration the friction-locking unit and/or the form-locking unit of the combined gear shift unit are disposed axially between the housing end wall on the drive side and the planetary gear stage associated with the combined gear shift unit, so that a particularly compact multi-stage transmission can be provided.\nIt is particularly advantageous if the form-locking unit of the combined gear shift unit is disposed axially between the friction-locking unit and the planetary gear stage associated with the combined gear shift unit, so that the form-locking unit can be configured in a particularly space-saving manner.\nIt is advantageous to place the form-locking unit on a small radius. Where appropriate the form-locking unit can be disposed on the same radius as the friction-locking element.\nFurthermore it is proposed that the multi-stage transmission comprises a further gear shift unit which is configured as a clutch and has a rotatable coupling element which is permanently connected non-rotatably to the transmission input shaft and an outer disc support which is configured at least partially integrally with an inner disc support of the friction-locking unit of the combined gear shift unit. In this way a space-saving design of the combined gear shift unit and the further gear shift unit can be realized. A “gear shift unit configured as a clutch” should in particular be understood to be a gear shift unit which is provided in order that the two rotatably disposed coupling elements thereof, which are rotatable independently of one another in an opened state, can be connected non-rotatably to one another in a closed state. A gear shift unit configured as a clutch is preferably disposed in a flux of force between at least two planetary gear stages. “Integrally” should in particular be understood to mean formed in one piece, such as for example by production from a casting. The outer disc support and the inner disc support configured integrally with the outer disc support are preferably provided in order to accommodate inner discs of the friction-locking unit of the combined gear shift unit and to accommodate outer discs of a friction-locking unit of the further gear shift unit.\nIn particular it is advantageous if the form-locking unit of the combined gear shift unit is disposed axially between the housing end wall on the drive side and the further gear shift unit, so that an alternative arrangement of the form-locking unit can be provided.\nFurthermore, it is advantageous if the combined gear shift unit has a first actuating unit provided for actuating the friction-locking unit as well as a second actuating unit provided for actuating the form-locking unit and the further gear shift unit has a third actuating unit, wherein the first actuating unit and the third actuating unit are disposed axially between the housing end wall on the drive side and the second actuating unit. In this way an advantageous actuation of the combined gear shift unit and the further gear shift unit can be achieved without the provision of additional installation space.\nIn an alternative embodiment it is advantageous if the combined gear shift unit has a first actuating unit provided for actuating the friction-locking unit as well as a second actuating unit provided for actuating the form-locking unit and the further gear shift unit has a third actuating unit, wherein the first actuating unit, the second actuating unit and the third actuating unit are disposed axially between the housing end wall on the drive side and the friction-locking unit of the combined gear shift unit. In this way the form-locking unit of the combined gear shift unit can be disposed radially inside the friction-locking unit of the combined gear shift unit and radially inside a friction-locking unit of the further gear shift unit, so that the circumferential speed of the form-locking tooth system can be further reduced. Furthermore, it may also be provided that closing directions of the actuating units are all aligned, so that a supply of an actuating medium to the actuating units for closing the friction-locking units and the form-locking unit can be simplified. A “closing direction” should in particular be understood to be a direction in which an actuating element of the actuating unit must be moved in order to close the friction-locking unit or the form-locking unit. In this context “oriented in the same direction” should in particular be understood to mean that the closing directions point in the same direction. An “actuating element” should in particular be understood to be an axially movable element which presses the inner and outer discs of a friction-locking unit against one another for closing the friction-locking unit or moves a form-locking tooth system of a form-locking unit into a corresponding form-locking tooth system of the form-locking unit for closing the form-locking unit.\nIn a further alternative embodiment it is advantageous if the combined gear shift unit has a first actuating unit provided for actuating the friction-locking unit as well as a second actuating unit provided for actuating the form-locking unit and the further gear shift unit has a third actuating unit, wherein the second actuating unit and the third actuating unit are disposed axially between the housing end wall on the drive side and the first actuating unit. In this way the form-locking unit of the combined gear shift unit can be disposed radially inside the friction-locking unit of the further gear shift unit in particular without affecting the installation space.\nMoreover, it is advantageous if, with respect to an axial succession starting from the internal combustion engine, the planetary gear stage associated with the combined gear shift unit is configured as a first planetary gear stage and a planetary gear stage adjacent to the planetary gear stage associated with the combined gear shift unit is configured as a second planetary gear stage, wherein the adjacent planetary gear stage has a transmission element which is configured as a planetary gear support and is permanently connected non-rotatably to a transmission element, configured as a ring gear, of the planetary gear stage associated with the combined gear shift unit, and has a transmission element which is configured as a ring gear and is permanently connected non-rotatably to a transmission element, configured as a sun gear, of a third planetary gear stage and to a transmission element, configured as a sun gear, of a fourth planetary gear stage and can also be connected non-rotatably to the transmission element, configured as a planetary gear support, of the planetary gear stage associated with the combined gear shift unit. In this way a particularly advantageous interconnection of the planetary gear stages can be realized. With respect to the axial succession of the planetary gear stages along an axial direction starting from the internal combustion engine, the first planetary gear stage is disposed in the first place, the second planetary gear stage is disposed in the second place, the third planetary gear stage is disposed in the third place and the fourth planetary gear stage is disposed in the fourth place. A “planetary gear stage adjacent to the planetary gear stage associated with the combined gear shift unit” should be understood to be a planetary gear stage which is disposed axially next in the sequence to the planetary gear stage associated with the combined gear shift unit.\nThe configuration of a multi-stage transmission described below constitutes a special exemplary embodiment, for which in principle there are further kinematically equivalent exemplary embodiments. For example in particular the planetary gear stages can be formed by means of single planetary gear sets, as is set out below. However, in a kinematically equivalent manner one of the planetary gear stages can also have a double planetary gear set, so that for the same mode of action only one stationary gear ratio of the planetary gear stage must be adapted. Moreover, a connection of a sun gear and a planetary gear support or of a ring gear and a planetary gear support can also be changed, without the kinematics of the multi-stage transmission changing.\nFurther advantages can be seen from the following description of the drawings. Four exemplary embodiments of the invention are shown in the drawings. The drawings, the description of the drawings and the claims contain numerous features in combination. Expediently, the person skilled in the art will also consider the features singly and combine them to form meaningful further combinations."}
{"text": "Generally, inflatable medical devices are introduced into closed circulatory system vertebrates and invertebrates in their uninflated form and are inflated when they reach the site of interest. Balloon catheters are examples of inflatable devices. A balloon catheter is a type of catheter which is bendable and has an inflatable balloon at its tip. These catheters may be used to enlarge a narrow opening or passage within the body. The catheter is advanced in a subject's blood vessel so that the balloon is in the correct position. The uninflated balloon is then inflated to perform the necessary procedure, and then deflated in order to be removed from the subject's blood vessel. When balloon catheters are inflated, they can occlude the opening or passage in which they are placed. Such an occlusion is generally temporary. Occasionally, it is necessary to inflate an element of a medical device that is not co-linear with the inflation lumen connecting the inflatable element to the operator. In such an arrangement, the inflatable element is set at an angle from the axis of the supporting catheter. When inflated, for the device to function properly, a relatively rigid, fixed orientation between the catheter body and the inflatable element must be achieved and maintained. In particular, this is true of medical devices which incorporate an inflatable toroid element in order to perform a particular function. Such toroids are typically fixed perpendicularly to the supporting catheter, and must be maintained in this orientation to function properly. Other devices with inflatable elements may also require that an inflatable element be disposed at an angle with relation to the supporting catheter. In each such instance, it is crucial that a structure exists to rigidly maintain this orientation. However, such a structure has not been described in the prior art.\nIt would therefore be desirable for an insertable, inflatable medical device to have a stable junction point through which an inflation fluid may pass and inflate an element which is not co-linear with the axis of inflation. Accordingly, there is a need in the art for such a structure that maintains a rigid orientation between an inflatable element and a catheter body, as one does not presently exist in the art."}
{"text": "The present invention relates to a system for recording and reproducing multiplex information. More particularly, the invention relates to a system for recording and reproducing multiplexed multichannel signals such as audio information in addition to video signals.\nVideo disks serving as information storage media have recorded thereon video signals and different information signals such as audio information signals. The video signals are frequency modulated. The audio signals are two-channel signals to allow for a stereophonic or bilingual program mode. The signals in each channel are recorded on the video disk by frequency modulating two independent audio carriers, together with the frequency-modulated video signals.\nThere is a demand for the multiplexing of other two-channel information, for example, in addition to the normal two-channel audio information. For reasons described later, there is no space available within the frequency spectrum recorded on a video disk for recording two additional channels. For multiplexing four-channel audio signals, furthermore, compatibility should be taken into account with existing recording and reproducing systems.\nIt is thus an object of the present invention to provide a multiplex information recording and reproducing method with which other recorded video information is not adversely affected, and which is compatible with existing recording and reproducing systems.\nAnother object of the present invention is to provide such a multiplex information recording and reproducing system having good reproducing characteristics."}
{"text": "1. Field of the Invention\nThe present invention relates to a multi-directional switch, and particularly relates to, for example, a multi-directional switch in which on-off operation of multiple contacts is dependent on the pressing directions of the switch.\n2. Description of the Prior Art\nFIG. 1 is a perspective view showing an example of a hand-held type game playing apparatus using a conventional character moving switch which constitutes the background of the present invention. Referring to FIG. 1, a body 1 comprises a liquid crystal display plate 2. The liquid crystal plate 2 comprises a plurality of characters formed with transparent electrodes so that a predetermined character is displayed according to development of a game or in response to operation of character moving switches 4a to 4d. The characters 3a, 3b, 3c and 3d illustrated in the drawing are selectively displayed by operation of character moving switches 4a, 4b, 4c and 4d. More particularly, the character 3a is displayed responsive to operation of the character moving switch 4a, the character 3b is displayed responsive to operation of the character moving switch 4b, the character 3c is displayed responsive to operation of the character moving switch 4c and the character 3d is displayed responsive to operation of the character moving switch 4d. Accordingly, by operating optionally the character moving switches 4a to 4d, a display is made as if characters 3a to 3d were moving. The body 1 further comprises game switches 5 and 6 for giving instructions to start a game of different difficulty and a current time switch 7 for giving instructions to indicate the current time.\nAs described above, the conventional character moving switches 4a to 4d are provided corresponding to the characters 3a to 3d, that is, corresponding to the moving directions of the characters. The character moving switches 4a and 4b are provided in the vicinity of the left end of the body 1, and the character moving switches 4c and 4d are provided in the vicinity of the right end of the body 1. Accordingly, if one plays a game with a hand-held type game playing apparatus in his hands, the character moving switches 4a and 4b must be operated with the left hand, and the character moving switches 4c and 4d must be operated with the right hand. Thus, it is very difficult to operate two switches with one hand and other two switches with the other hand. As a result, one must take a great care in operation of the character moving switches 4a and 4d, which makes the game dull. On the other hand, if a larger number of character moving switches are provided to increase character moving directions, more variety will be given to the game, so that one will not easily get bored with it. Thus, in using conventional character moving switches, there is a disadvantage that either variety of a game or simplicity in operation must be disregarded.\nIn order to eliminate the above described disadvantage, the character moving switches 4a to 4d might be disposed collectively in a certain position in the body 1 so that one can operate them with one hand. However, such collective disposition will encounter another problem in that two or more than two character moving switches are often pressed simultaneously. Conversely, this means that the character moving switches 4a to 4d must be separated from each other with a certain distance. Accordingly, space for positioning the character moving switches 4a to 4d should be made larger, and consequently sufficient space cannot be provided for display of the liquid crystal display plate 2. As a result, the content of a game organized in the liquid crystal display plate 2 has to be limited.\nIn the foregoing, the disadvantage of character moving switches used in a hand-held type game playing apparatus was explained. However, the same disadvantages or problems as described above are also found in various conventional apparatus where control for selection of modes or moving directions is made by means of switches, if switches are provided corresponding to the modes or moving directions."}
{"text": "In the traditional client-server model, the majority of processing and data storage occurs at the client. In the emerging Software as a Service (“SaaS”) market, the majority of the processing occurs at the server and the data is centralized to take advantage of the connectivity the Internet provides. The power of the connectivity and information sharing is leading to massively scalable applications that support hundreds of thousands, up to hundreds of millions of users.\nMassively scalable applications are creating many new challenges in managing user loads and storage systems in an automated fashion. One of theses challenges is the ability to accurately predict when capacity will be needed in data-heavy applications, such as email, file storage, and online back-up, and also in non-data heavy applications. Making this prediction is difficult because the limitations which can affect available overall load take many forms, including utilization of processor, memory, I/O load (comprising reads per second, writes per second, total transactions per second, and number of ports being utilized), network space, disk space, an application or applications, and power, and these forms are continually changing. Additionally, making this prediction is difficult because it is hard for an administrator to distinguish between a short-term load that can be reached due to a burst of activity and long-term growth that requires increased capacity. Administrators also do not have useful numbers on the different types of utilization of the server systems. As a result, administrators are often in a reactive mode when trying to address capacity concerns. Further, administrators do not have any methods for accurately assessing costs, such as the cost per user per service especially when multiple servers are required for a single service or for assessing whether to maximize the configuration of a server or purchase a new server to add capacity."}
{"text": "In recent years, along with a development of multimedia applications such as picture, audio and text, it has become general to handle all sorts of media in an integrated way. However, an information compression technique for data is dispensable for its storage and transmission since a digitalized picture contains an enormous amount of data. On the other hand, a standardization of compression techniques is also important for interoperating compressed picture data. The standards of picture compression techniques include H.261, H.263 established by the ITU (International Telecommunication Union) and MPEG (Moving Picture Experts Group)-1, MPEG-2 and MPEG-4 established by the ISO (International Organization for Standardization).\nAn inter-picture prediction which accompanies motion compensation can be cited as a technique shared among these moving picture coding methods. In the motion compensation based on these moving picture coding methods, a picture of an input image is divided into blocks, each of which has a predetermined size, and a predictive image is generated for each block using motion vectors, respectively indicating a motion between pictures. The following predictions are employed for the inter-picture prediction according to the MPEG: a forward prediction for a prediction using a single picture whose display time is earlier than that of a current picture to be coded; a backward prediction for a prediction using a single picture whose display time is later than that of the current picture; a bi-directional prediction for a prediction using two pictures, that is, one picture whose display time is earlier than the current picture and the other picture whose display time is later than that of the current picture (see reference, for example, ISO/IEC 14496-2:1999(E) Information technology—coding of audio-visual objects Part 2: Visual (1999-12-01) pp 150 7.6.7 Temporal prediction structure).\nIn the MPEG, a reference picture to be used is determined uniquely depending on the type of inter-picture prediction and an arbitrary reference picture cannot be selected. In the meantime, a bidirectional prediction which is expanded so that two arbitrary reference pictures can be selected out of a plurality of coded pictures stored in a picture memory regardless of the display time of the current picture is taken under the consideration in the H.264 which is presently under the process of standardization by the ITU.\nFIG. 1 is a block diagram showing a structure of a moving picture coding apparatus according to the H.264. The conventional moving picture coding apparatus shown in FIG. 1 is an apparatus for executing a moving picture coding method which allows a selection of two arbitrary reference pictures from plural coded pictures when the inter-picture prediction is operated.\nThis moving picture coding apparatus includes, as shown in FIG. 1, a motion estimation unit 301, a pixel interpolation unit 102, a subtractor 103, a picture coding unit 104, a picture decoding unit 105, an adder 106, a variable length coding unit 302, a multi-picture buffer 108 and a switch 109.\nThe moving picture coding apparatus divides an inputted image data Img into blocks and performs processing for each of the blocks. The subtractor 103 subtracts a predictive image data Pred from the image data Img inputted to the moving picture coding apparatus and outputs it as residual data Res. The picture coding unit 104 performs picture coding processing such as orthogonal transformation and quantization on the inputted residual data Res and outputs it as coded residual data ERes including quantized orthogonal transformed coefficients. The picture decoding unit 105 performs picture decoding processing such as inverse quantization and inverse orthogonal transformation on the inputted coded residual data ERes and outputs it as decoded residual data DRes. The adder 106 adds the decoded residual data DRes to the predictive image data Pred and outputs it as reconstructed image data Recon. Out of the reconstructed image data Recon, the data having the possibility to be used for reference in the subsequent inter-picture prediction is stored in the multi-picture buffer 108.\nHere, an interpolation prediction using two reference pictures performed by the conventional moving picture coding apparatus is described with reference to FIG. 2. FIG. 2 is a conceptual diagram of the interpolation prediction using plural reference pictures. Here, a picture Pic is a current picture to be coded. Pictures FwRef1˜FwRef3 represent coded pictures respectively having a display time earlier than that of the current picture whereas pictures BwRef1˜BwRef3 represent coded pictures respectively having a display time later than that of the current picture. A block Blk1 is predicted using pixel values in a reference block RefBlk11 included in the picture FwRef3 whose display time is earlier than that of the current picture Pic and pixel values in a reference block RefBlk12 included in the picture BwRef1 whose display time is later than that of the current picture Pic. A block Blk2 is predicted using pixel values in reference blocks RefBlk21 and RefBlk22 included in two pictures FwRef1 and FwRef2 respectively having a display time earlier than that of the current picture. A block Blk3 is predicted using pixel values in reference blocks RefBlk31 and RefBlk32 included in two pictures BwRef1 and BwRef2 respectively having a display time later than that of the current picture. Namely, a result of interpolating pixels in the areas corresponding to the two reference blocks using a predetermined method such as the one using an average value is considered to be a predictive image. The characteristics of the conventional moving picture coding apparatus is to perform prediction on a block-by-block basis using arbitrary two reference pictures as shown in FIG. 2. A method for predicting with the use of two arbitrary reference pictures as described above is called “plural reference picture interpolation prediction” hereinafter. The prediction method includes a method in which a block included in a single arbitrary picture is used directly as a predictive image and the intra-picture prediction other than the method of generating a predictive image using the pixel interpolation as described above, and it is possible to switch the prediction method on a block-by-block basis.\nThe motion estimation unit 301 determines a prediction type for the block, reference pictures and motion vectors to be used for inter-picture prediction performed on the inputted current block to be coded and outputs a prediction type PredType, reference picture numbers RefNo1, RefNo2, and motion vectors MV1, MV2. The motion estimation 301 outputs two picture numbers and two motion vectors since two reference pictures are selected when plural reference picture interpolation prediction is operated. Here, the multi-picture buffer 108 outputs a reference block RefBlk1 corresponding to the reference picture number RefNo1 and the motion vector MV1 and a reference block RefBlk2 corresponding to the reference picture number RefNo2 and the motion vector MV2. The pixel interpolation unit 102 performs interpolation for the pixels with respect to the two reference blocks RefBlk1 and RefBlk2 using average value and outputs it as-an interpolated block RefPol. On the other hand, in the case of using an inter-picture prediction other than a plural reference picture interpolation prediction, the motion estimation unit 301 selects a single reference picture, and therefore, outputs a single reference picture number RefNo1 and a single motion vector MV1. In this case, the multi-picture buffer 108 outputs a reference block RefBlk with respect to the reference picture number RefNo1 and the motion vector MV1.\nWhen the prediction type determined by the motion estimation unit 301 indicates a plural reference picture interpolation prediction, the switch 109 is switched to a “1” side and the interpolated block RefPol is used as a predictive image data Pred. When the prediction type PredType indicates an inter-picture prediction other than a plural reference picture interpolation prediction, the switch SW11 is switched to a “0” side and the reference block RefBlk is used as a predictive image data Pred. The variable length coding unit 302 performs variable length coding on the coded residual data ERes, the prediction type PredType, the reference picture numbers RefNo1, RefNo2 and the motion vectors MV1, MV2 and then outputs them as coded moving picture data Str0.\nFIG. 3 is a conceptual diagram of a data format of coded moving picture used by the conventional moving picture coding apparatus. Coded data equivalent to a single picture, Picture, is composed of coded data equivalent to a single block, Block, where each block composes a picture, and the like. Here, the coded data equivalent to a single block, Block, presents coded data of a block on which a plural reference picture interpolation prediction is performed, and includes the reference picture numbers RefNo1, RefNo2, the motion vectors MV1, MV2, with respect to the two reference pictures, the prediction mode PredType, and the like, in the coded data.\nFIG. 4 is a block diagram showing a structure of the conventional moving picture decoding apparatus. The moving picture decoding apparatus includes, as shown in FIG. 4, a variable length decoding unit 601, a motion compensation unit 602, a picture decoding unit 404, an adder 405, a pixel interpolation unit 406, a multi-picture buffer 407 and a switch 408.\nThe variable length decoding unit 601 performs variable length decoding on the inputted coded image data Str0 and outputs the coded residual data ERes, the motion vectors MV1, MV2, the reference picture numbers RefNo1, RefNo2 and the prediction type PreType. The picture decoding unit 404 performs picture decoding processing such as inverse quantization and inverse orthogonal transformation on the inputted coded residual data ERes and outputs decoded residual data DRes. The adder 405 adds the decoded residual data DRes to the predictive image data Pred and outputs it as decoded image data DImg outside the moving picture decoding apparatus. The multi-picture buffer 407 stores the decoded image data DImg for inter-picture prediction.\nThe motion compensation unit 602 outputs reference picture numbers NRefNo1, NRefNo2 of the reference blocks necessary for inter-picture prediction according to the prediction type PredType as well as the motion vectors MV1, MV2 and instructs the multi-picture buffer 407 to output the reference blocks. When the prediction type PredType indicates a plural reference picture interpolation prediction, the multi-picture buffer 407 outputs the reference block RefBlk1 corresponding to the reference picture number NRefNo1 and the motion vector NMV1 as well as the reference block RefBlk2 corresponding to the reference picture number NRefNo2 and the motion vector NMV2. The pixel interpolation unit 406 interpolates the pixels in the two reference blocks RefBlk1 and RefBlk2 using the average value. On the other hand, when the prediction type PredType indicates an inter-picture prediction method other than a plural reference picture interpolation prediction, the multi-picture buffer 407 outputs the reference block RefBlk corresponding to the reference picture number NRefNo1 and the motion vector NMV1.\nWhen the prediction type PreType indicates a plural reference picture interpolation prediction, the switch 408 is switched to a “0” side and an interpolated block RefPol is used as a predictive image data Pred. Thus, the moving picture decoding apparatus decodes the coded moving picture data Str0 through the processing described above and outputs it as decoded image data DImg.\nMeanwhile, under the moving picture coding method based on the MPEG-4, a plural reference picture interpolation prediction method called “direct mode” is defined for a picture type, called “bidirectional predictive picture”, employing a plural reference picture interpolation prediction. It is defined as a method to abbreviate the motion vectors and the reference picture numbers included in the coded data of the block by calculating the motion vectors with respect to two reference pictures used for the generation of the predictive image by means of interpolation using the coded motion vectors.\nFIG. 5 is an illustration for a case of using the direct mode defined in the MPEG-4. Here, a picture Pic represents a current picture to be coded, a picture Ref1 represents a reference picture whose display time is earlier than that of the current picture Pic and a picture Ref2 represents a reference picture whose display time is later than that of the current picture Pic whereas a block Blk represents a current block to be coded and a block Blk0 represents a block whose position is same as that of the current block Blk in the reference picture Ref2. A motion vector MV01 represents a forward reference motion vector using the picture Ref1 as a reference picture for coding the block Blk0, a motion vector MV1 represents a motion vector of the current block with respect to the reference picture Ref1, a motion vector MV2 represents a motion vector of the current block with respect to the reference picture Ref2, a block RefBlk1 represents a reference block to be referred to by the motion vector MV1 and a block RefBlk2 represents a reference block to be referred to by the motion vector MV2.\nAs for the two pictures to be used for reference by the current block Blk, the picture Ref2 whose display time is later than and is closest to the current picture is used as a backward reference picture whereas the picture Ref1, which has been used for reference by the block Blk0 at the time of coding, is used as a forward reference picture.\nFor the calculation of the motion vectors, it is assumed that either the motion is constant or no motions are found in comparing the pictures. Here, with an assumption that a differential value between the display time of the current picture Pic and that of the reference picture Ref1 is TRD1, a differential value between the display time of the reference picture Ref1 and that of the reference picture Ref2 is TRD2, and a differential value between the display time of the current picture Pic and that of the reference picture Ref2 is TRD3, the motion vectors MV1 and MV2 to be used for coding the current block can be calculated respectively using the following equations:MV1=MV01×(TRD1/TRD2)  (Equation A)MV2=−MV01×(TRD3/TRD2)  (Equation B)\nUsing the above method, the reference pictures and the motion vectors in the case of using a direct mode can be determined. The processing in the case of using a direct mode as described above, performed by the moving picture coding apparatus, is executed by the motion estimation unit 301 shown in the block diagram illustrating the conventional moving picture coding apparatus in FIG. 1. The processing for the case of using a direct mode described above, performed by the moving picture decoding apparatus, is executed by the motion compensation unit 602 shown in the block diagram illustrating the conventional moving picture decoding apparatus in FIG. 4.\nWhen a moving picture, in which a motion between the pictures is small, is inter-picture coded, a predictive error between the pictures become very small and most of the coded residual data ERes become “0” by performing picture coding processing such as quantization. A case in which the entire coded residual data ERes resulted from the inter-picture prediction using the reference pictures and the motion vectors of the current block is “0” in the coding in which the motion vectors and the reference pictures are determined using a predetermined method without coding them, as in the case of using a direct mode as described above, is defined as one of the prediction types PredType called “skip mode”. In using a skip mode, only the prediction type PredType indicating the skip mode is transmitted, therefore, coding of a block requires a very small code amount. The efficiency of coding can be further improved by assigning variable length code that is shorter than other prediction types to this skip mode or by run-length coding the number of consecutive blocks used for the skip mode.\nIn the H.264 described above, “skip mode” is defined as a case in which the entire coded residual data equivalent to a single block obtained by the inter-picture prediction using a direct mode is assumed to be “0”. The following processing is performed when a block is coded using a skip mode by the moving picture coding apparatus shown in FIG. 1. The motion estimation unit 301 outputs the reference picture numbers RefNo1, RefNo2, the motion vectors MV1, MV2 as well as the prediction type PredType indicating a skip mode. The variable length coding unit 302 performs variable length coding only for the prediction type PredType and outputs it as coded moving picture data Str0 through the processing explained above, when the prediction type PredType indicates a skip mode. The following processing is performed when the coded data of the block coded using a skip mode is inputted to the moving picture decoding apparatus shown in FIG. 4. The variable length decoding unit 601 performs variable length decoding on the prediction type PredType. When the prediction type PredType indicates a skip mode, the motion compensation unit 602 outputs, through the processing operated in the case of direct mode explained above, the reference picture numbers NRefNo1, NRefNo2, the motion vectors NMV1, NMV2 as well as the prediction type PredType indicating a skip mode.\nIn the H.264 as described above, arbitrary reference pictures can be selected out of a plurality of coded pictures regardless of the display time of the current picture. However, the arbitrary reference pictures are selected by performing motion estimation for the plurality of the coded pictures in this case, therefore, the processing burden caused by the motion estimation becomes very large. The plural reference picture interpolation prediction also contains a problem of degrading the coding efficiency since it requires coding of reference picture numbers and motion vectors for every two reference pictures.\nFurthermore, when inter-picture prediction is performed for a picture using a picture whose display time is later than that of the current picture as a reference picture, as in the case of bidirectional prediction described in the conventional technique, the picture has to be coded in an order different from a display order, which causes a delay. In a case of real time communication such as a videophone, bi-directional predictive pictures cannot be used because of the delay. In the H.264, however, two arbitrary reference pictures can be selected regardless of display order information, therefore, the delay caused by coding can be eliminated by performing a plural reference picture interpolation prediction with a selection of two pictures respectively having a display time which is earlier than that of the current picture. However, the picture whose display time is later than that of the current picture is not stored in the multi-picture buffer, therefore, the direct mode conventionally used for determining the motion vectors using the picture whose display time is later than that of the current picture as described above cannot be employed."}
{"text": "1. Field of the Invention\nThe present invention relates to integrated circuit lead frames and methods of production thereof. In particular, this invention relates to plastic lead frames with a conductive coating or material contained therein used for packaging integrated circuits and methods of manufacturing the same.\n2. State of the Art\nIntegrated circuit (IC) chips are enclosed in plastic packages that provide protection from hostile environments and enable electrical interconnection to printed-circuit boards. During a manufacturing process, the IC chip is typically attached to a die paddle of a conventional lead frame or suspended from the lead fingers of a leads-over-chip (LOC) lead frame using an adhesive such as epoxy or double-sided tape, and subsequently encapsulated with a dense and rigid plastic by a transfer molding process. In essence, the lead frame forms the backbone of the molded plastic IC package.\nLead frames typically perform many functions such as: (1) a holding fixture that indexes with tool-transfer mechanisms as the package proceeds through various assembly operations, (2) a dam that prevents plastic from rushing out between leads during the molding operation, (3) a chip attach substrate, (4) a support matrix for the plastic, and (5) an electrical and thermal conductor from chip to board.\nTraditionally, lead frames are fabricated from a strip of sheet metal by stamping or chemical milling operations. There are many different metal alloy compositions which are commercially available for producing lead frames. For example, Rao R. Tummala and Eugene J. Rymaszewski, \"Microelectronics Packaging Handbook,\" Table 8-4, 1989, provide 16 different alloys available from 9 different manufacturers. Lead frame material selection depends on many factors such as cost, ease of fabrication, strength, thermal conductivity, and matched coefficient of thermal expansion (CTE). A close match of CTE between the silicon die and the lead frame is required to avoid chip fracture from differential expansion rates.\nThe most widely used metal for lead frame fabrication is Alloy 42 (42% Nickel-58% Iron). Alloy 42 has a CTE near silicon and good tensile strength properties. The disadvantage of Alloy 42 is that it has low thermal conductivity. Since the lead frame is the main conduit by which heat flows from the chip to the environment and printed circuit board, this can have a profound effect on the package thermal resistance after prolonged device operation.\nA layered composite strip, such as copper-clad stainless steel, was developed to emulate the mechanical properties of Alloy 42 while increasing thermal conductivity. However, copper-clad stainless steel is somewhat more expensive to manufacture than Alloy 42. When manufacturing copper-clad stainless steel lead frames, the cladding is accomplished by high-pressure rolling of copper foil onto a stainless steel strip, followed by annealing the composite to form a solid-solution weld. While copper alloys provide good thermal conductivity and have a CTE near that of low-stress molding compounds, there is a substantial CTE mismatch with respect to silicon.\nWhile numerous alloys have been developed to solve problems with thermal conductivity, CTE mismatch, and strength, other important factors such as ease of fabrication and cost have not improved as readily.\nConventional methods for making lead frames for integrated circuit devices are described in U.S. Pat. No. 3,440,027. The use of a plastic support structure in a method of forming metal lead frames is described in U.S. Pat. No. 4,089,733 (hereinafter the \"'733\" patent). The plastic support structure of the '733 patent solves the problem of deformed and misaligned lead fingers resulting from stress during the bonding process by supporting the lead fingers with a plastic structure. However, the '733 patent requires a metal lead frame in addition to the plastic support structure with its attendant costs. A method of manufacturing multi-layer metal lead frames is disclosed in U.S. Pat. No. 5,231,756 (hereinafter the \"'756\" patent). The '756 patent provides an improvement in aligning power and ground planes for use in a multi-layer lead frame where such planes are necessary. However, the number of steps required to manufacture such multi-layer lead frames will not solve the problem of decreasing costs. In short, none of these related art appear to disclose methods of producing low-cost lead frames made from materials not structurally based on metal.\nSince packaged ICs are produced in high volumes, a small decrease in the cost per packaged IC can result in substantial savings overall. Accordingly, there is a need in the industry for a low-cost plastic lead frame with suitable characteristics for IC packaging."}
{"text": "In a storage device using a semiconductor memory or the like, error correction processing is performed in data decoding processing. There is a chien search as a method of specifying error locations of data in the error correction processing. In the chien search, exclusive-OR operations multiplying α to the power, which multiply coefficients of an error location polynomial (t-order or less) except for a constant term by α, α2, . . . , αt, and exclusive-OR operations, which obtain the sum of the respective results of the operations on a Galois field, are performed to specify error locations by searching roots. Since error locations are determined in every cycle in a chien search circuit, all the exclusive-OR operations are performed in one cycle. Meanwhile, it is possible to improve error correction capability by increasing the order of the Galois field or the order of the error location polynomial in the chien search circuit, but the number of stages of the exclusive-OR operations is increased."}
{"text": "Universal Serial Bus (USB) is an industry standard developed that defines the cables, connectors and communications protocols used in a bus for connection, communication, and power supply between electronic devices.\nIn many portable products, USB is often used both as a communication port and as a power delivery port to accommodate battery charging. For example, a standard USB 2.0 compliant port may provide a maximum power delivery of 7.5 W (5V at 1.5 A) to a dedicated charging port that may be used to recharge the battery of a portable device. However, as the battery capacities of portable devices are increasing, for example, from 5600 mAh to 8000 mAh and 10000 mAh, the charging time for these devices increases accordingly. For example, using a standard USB 2.0 compliant port, it takes about 2 hours and 40 minutes to recharge a 5600 mAh battery, but it takes 4 hours and 45 minutes to recharge a 10000 mAh.\nFaster charging times may be achieved by increasing the charging current and/or charging voltage. USB Power Delivery (PD) specification specifies using certified PD aware USB cables with standard USB type A/B connectors to deliver up to 100 W of power at 20 V. USB PD requires the use of PD-aware cables. For example, PD-aware cables with USB-micro B/AB connectors may support a maximum power of up to 60 W at 20 V, 36 W at 12 V and 10 W at 5 V. Consequently, under new USB standards, currents up to 5 A and power up to 100 W may be transferred through power delivery aware cables.\nUSB specification requires specification on the electrical parameters for the USB cables. For example, USB specification requires that USB cables have an American Wire Gauge (AWG) rating of 24-26, and further require that micro-B connectors have less than 50 mΩs.\nHowever, USB cables including PD-aware cables may not be able to safely handle the large currents. For example, the cables may become defective during use due to field stress. Additionally, some of the USB manufacturers may not have good quality control. The large power supplied through a defective USB can pose a significant safety risk. Similarly, due to backward compatibility, end users may not be aware of the necessity to use the newer PD-aware cables."}
{"text": "Cellular wireless networks commonly provide both voice and data capability. A data communications layer is typically overlaid onto the voice communications layer, while utilizing the same network. For example, cellular devices utilizing the global system for mobile communications (GSM) voice standard are able to use the general packet radio service (GPRS) layer of the network for data communications. In these systems, a cellular device must register separately for the voice and the data networks. Protocols for network selection and registration were designed for human users primarily using voice communications, and in some instances these protocols require manual reconfiguration if there is a problem with or a different selection is desired for the data network. However, a problem arises in the event that a manual configuration, correction, or other input is required to correct a problem with the data network and there is no one available to provide the manual configuration, correction, or other input."}
{"text": "With the advent of the Internet, email has become prevalent in digital communications. For example, email messages are exchanged on a daily basis to conduct business, to maintain personal contacts, to send and receive files, etc. Unfortunately, undesired email messages have also become prevalent with increased email traffic. Often, these email messages are unsolicited advertisements, which are often referred to as “junk mail” or “spam.” Currently, software applications exist, which remove some of the spam or junk mail from a recipient's email account, thereby reducing clutter in the recipient's email account. Some of these applications remove email messages that contain a particular text or content (e.g., large image files, etc.) that may indicate that the email message is spam or junk mail. Email messages that are determined to be spam or junk mail are then either removed (e.g., permanently deleted) or stored in a designated folder (e.g., “trash” folder, “junk mail” folder, etc.). Unfortunately, some of the algorithms used to detect spam or junk mail may be quite complicated and cumbersome. Due to the complexity of current anti-spam or antijunk-mail algorithms, a need exists in the industry for simpler solutions."}
{"text": "Modern high volume package delivery systems often include package conveying systems that accept packages from one or more loading stations, and transports the packages to a variety of output destinations such as chutes, bins, and subsequent conveyor systems.\nOne of the most conventional types of conveyors is a belt conveyor, which includes the use of an endless flexible belt which passes over at least two cylindrical rollers, one of which is a drive roller. Packages are placed atop the upwardly-directed \"working\" surface of the belt conveyor, and are transported in a generally straight direction from end of the conveyor to the other.\nAlthough such conventional belt conveyors have advantages, including simplicity, they also have disadvantages in that they conventionally only provide a \"straight-line\" transporting feature. This can be restrictive to package delivery system designers who may have the need to move a package or other product from an origin through a tortuous, curved, path to a destination.\nTherefore, it has been known in the prior art to provide flexible conveyor chains such as shown in U.S. Pat. No. 3,776,349 to Kampfer, entitled \"Fabricated Conveyor Chain\", which discloses the concept of providing a fabricated flexible conveyor chain, which includes a plurality of link units 11 linked together by a plurality of pivot pins 12. Although the pivot pins 12 provide a linking feature between the link units, they fit loosely enough within their mounting holes to allow sideward relative pivoting of the link units. A similar type of \"hard pin\" connection is disclosed in U.S. Pat. No. 3,262,550 to Kampfer, entitled \"Conveyor Chain\", in U.S. Pat. No. 2,884,118 to Williams, entitled \"Articulator Conveyor Chain\", and also in U.S. Pat. No. 5,176,247 to Counter et. al., entitled \"Sideflexing Conveyor Chain including Low Centerline Hinge Pin\".\nAlthough such \"hard pin\" connection configurations as described above have their advantages, they have disadvantages in that they tend to be complex, expensive, noisy, and difficult to maintain. Furthermore, they tend to provide a multitude of hard \"pinch points\", which are disadvantageous when in a human workplace environment. A \"smooth\" but flexible conveyor configuration is disclosed in U.S. Pat. No. 4,084,687 to Lapeyre, entitled \"Conveyor Having Resilient Conveying Surface\", but this configuration appears to be quite complex, requiring the use of link members 20 which are linked to modules 20 to support and convey flexible members 32 supported thereon.\nTherefore, there is a need in the art to provide a package conveyor system which can transport packages or other items along a tortuous path, yet is simple in construction, quiet in operation, and cost-effective to manufacture, operate, and maintain."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrode of a plasma display panel (PDP). More particularly, the present invention relates to a technology for improving exposure latitude and imparting improved PDP quality by modifying a photo-patternable paste.\n2. Technical Background\nIn PDP's, 1000 or more electrodes measuring 80 to 100 μm in width by 0.5 to 1.5 m in length are formed. These electrodes are preferred to be free of defects and have superior linearity. Silver is typically used for the conductive metal particles contained in the electrodes.\nPhotolithography has recently become widely used to form the electrodes in PDP. In this method, a photo-patternable paste is first coated. A pattern is then formed by exposing and developing a film composed of the photo-patternable paste. Accompanying the increased size of glass substrates in recent years, proximity exposure has come to be used during exposure. Proximity exposure refers to a system whereby exposure is carried out while providing a gap of several hundred micrometers between a photomask and the coated film to be exposed.\nIn the case of pattern formation using a photo-patternable paste, it is desirable to increase the process latitude. It is desirable to increase the degree of freedom of the exposure conditions. In the case of low process latitude, the pattern shape may change easily depending on the exposure conditions, or the width of the electrodes may become uneven. Photosensitivity to changes in the pattern shape depending on exposure conditions signifies a low degree of freedom, and may even lead to restrictions on the production process. When carrying out proximity exposure, in particular, the width of the electrode formed tends to be greater than the designed electrode width, resulting in the decrease of sharpness of image. Measures to avoid this decrease in sharpness are needed.\nJapanese Patent Application Laid-open No. 2003-195487 discloses a technology in which an ultraviolet absorber is added for the purpose of forming a fine, sharp transparent electrode. According to this publication, addition of the ultraviolet absorber makes it possible to diminish scattering of exposure light within the paste. Examples of ultraviolet absorbers include azo-based dyes, amino ketone-based dyes, xanthene-based dyes, quinoline-based dyes, anthraquinone-based dyes, benzophenone-based dyes, diphenylcyanoacrylate-based dyes, triazine-based dyes and p-aminobenzoic acid-based dyes.\nHowever, if a dye-based ultraviolet absorber like that disclosed in Japanese Patent Application Laid-open No. 2003-195487 is used, since the ultraviolet absorber itself absorbs light, the paste that is produced is also colored. In a process for producing a PDP using a conductive paste, defect inspections are frequently carried out using visible light to confirm the coated state of the paste. At this time, if a dye which absorbs visible light is contained in the paste, there is the possibility of locations where defects are not present being detected. In addition, if the absorbance of visible light changes as a result of adding an ultraviolet absorber, there is the risk of having a detrimental effect on the efficiency of the photoinitiator. For example, since some photoinitiators have an absorbance range in the vicinity of up to 450 nm, in the case of using such a photoinitiator, a change in the absorbance in that range also ends up having an effect on the reflection efficiency of the photoinitiator."}
{"text": "The present invention relates to image-forming optical systems. More particularly, the present invention relates to a decentered optical system with a reflecting surface having a power for use in optical apparatus using digital still cameras, film scanners, and endoscopes. a small-sized image pickup device, e.g. video cameras,\nRecently, with the achievement of small-sized image pickup devices, image-forming optical systems for use in video cameras, digital still cameras, film scanners, endoscopes, etc. have also been demanded to be reduced in size and weight and also in cost.\nIn the general rotationally symmetric coaxial optical systems, however, optical elements are arranged in the direction of the optical axis. Therefore, there is a limit to the reduction in thickness of the optical systems. At the same time, the number of lens elements unavoidably increases because it is necessary to correct chromatic aberration produced by a rotationally symmetric refracting lens used in the optical systems. Therefore, it is difficult to reduce the cost in the present state of the art. Under these circumstances, there have recently been proposed optical systems designed to be compact in size by giving a power to a reflecting surface, which produces no chromatic aberration, and folding an optical path in the optical axis direction.\nJapanese Patent Application Unexamined Publication (KOKAI) Number [hereinafter referred to as “JP(A)”] 7-333505 proposes to reduce the thickness of an optical system by giving a power to a decentered reflecting surface and thus folding an optical path. In an example thereof, however, the number of constituent optical members is as large as five, and actual optical performance is unclear. No mention is made of the configuration of the reflecting surface.\nJP(A) 8-292371, 9-5650 and 9-90229 each disclose an optical system in which an optical path is folded by a single prism or a plurality of mirrors integrated into a single block, and an image is relayed in the optical system to form a final image. In these conventional examples, however, the number of reflections increases because the image is relayed. Accordingly, surface accuracy errors and decentration accuracy errors are transferred while being added up. Consequently, the accuracy required for each surface becomes tight, causing the cost to increase unfavorably. The relay of the image also causes the overall volumetric capacity of the optical system to increase unfavorably.\nJP(A) 9-222563 discloses an example of an optical system that uses a plurality of prisms. However, because the optical system is arranged to relay an image, the cost increases and the optical system becomes large in size unfavorably for the same reasons as stated above.\nJP(A) 9-211331 discloses an example of an optical system in which an optical path is folded by using a single prism to achieve a reduction in size of the optical system. However, the optical system is not satisfactorily corrected for aberrations.\nJP(A) 8-292368, 8-292372, 9-222561, 9-258105 and 9-258106 all disclose examples of zoom lens systems. In these examples, however, the number of reflections is undesirably large because an image is relayed in a prism. Therefore, surface accuracy errors and decentration accuracy errors of reflecting surfaces are transferred while being added up, unfavorably. At the same time, the overall size of the optical system unavoidably increases, unfavorably.\nJP(A) 10-20196 discloses an example of a two-unit zoom lens system having a positive front unit and a negative rear unit, in which the positive front unit comprises a prism of negative power placed on the object side of a stop and a prism of positive power placed on the image side of the stop. JP(A) 10-20196 also discloses an example in which the positive front unit, which comprises a prism of negative power and a prism of positive power, is divided into two to form a three-unit zoom lens system having a negative unit, a positive unit and a negative unit. However, the prisms used in these examples each have two transmitting surfaces and two reflecting surfaces, which are all independent surfaces. Therefore, a relatively wide space must be ensured for the prisms. In addition, the image plane is large in size in conformity to the Leica size film format. Accordingly, the prisms themselves become unavoidably large in size. Furthermore, because the disclosed zoom lens systems are not telecentric on the image side, it is difficult to apply them to image pickup devices such as CCDs. In either of the examples of zoom lens systems, zooming is performed by moving the prisms. Accordingly, the decentration accuracy required for the reflecting surfaces becomes tight in order to maintain the required performance over the entire zooming range, resulting in an increase in the cost.\nWhen a general refracting optical system is used to obtain a desired refracting power, chromatic aberration occurs at an interface surface thereof according to chromatic dispersion characteristics of an optical element. To correct the chromatic aberration and also correct other ray aberrations, the refracting optical system needs a large number of constituent elements, causing the cost to increase. In addition, because the optical path extends straight along the optical axis, the entire optical system undesirably lengthens in the direction of the optical axis, resulting in an unfavorably large-sized image pickup apparatus.\nIn decentered optical systems such as those described above in regard to the prior-art, an imaged figure or the like is undesirably distorted and the correct shape cannot be reproduced unless the formed image is favorably corrected for aberrations, particularly rotationally asymmetric distortion.\nFurthermore, in a case where a reflecting surface is used in a decentered optical system, the sensitivity to decentration errors of the reflecting surface is twice as high as that in the case of a refracting surface, and as the number of reflections increases, decentration errors that are transferred while being added up increase correspondingly. Consequently, manufacturing accuracy and assembly accuracy, e.g. surface accuracy and decentration accuracy, required for reflecting surfaces become even more strict."}
{"text": "The present invention relates to the activation of fluid bed oxidation catalysts. More particularly the present invention relates to the activation of fluid bed oxidation catalysts useful in the preparation of maleic anhydride from 4-carbon atom hydrocarbons, including n-butane.\nOxidation catalysts containing the mixed oxides of vanadium and phosphorus have been utilized to produce maleic anhydride from 4-carbon atom hydrocarbons, such as n-butane. Methods have been investigated for activating, or increasing the catalytic activity of these catalysts. It is taught in the literature to \"condition\" vanadium phosphate-containing maleic anhydride catalysts under the flow of a low level of hydrocarbon in air, such as 0.2 volume percent to 2 volume percent hydrocarbon in air at temperatures of 300.degree. C. to 600.degree. C., as in U.S. Pat. No. 4,171,316.\nU.S. Pat. No. 3,985,775 discloses the preparation of vanadium phosphorus mixed oxide catalysts prepared by a method requiring heating the catalyst precursor at about 350.degree. C. to 410.degree. C. to drive off at least a portion of hydration, and then at a higher temperature. This method is taught as suitable for use in fixed bed operation.\nAttempting to \"condition\" a fluid bed catalyst with low levels of hydrocarbon in air under normal operating conditions has been found to have little beneficial effect.\nActivation procedures for fixed bed catalysts are not readily applicable for use with fluid bed catalysts. In fixed bed operation, a characteristic exotherm or hot spot is generated in the catalyst bed within the reactor vessel. A substantial portion of the total reaction takes place at this hot spot. The location of the hot spot is dependent upon the location of the oxygen and feed stock inlets, the dimensions of the vessel, and operating conditions such as flow rates, temperature and pressure. Fixed bed activation procedures generally may concentrate in the hot spot of the catalyst bed, having little effect on the remainder of the catalyst in the reactor.\nIn fluid bed operation, such a hot spot is not generated. Catalyst particles are not localized, but move throughout the reactor vessel, and substantially the entire portion of the catalyst within the reactor vessel contacts the feed to contribute to the reaction."}
{"text": "A critical response of the innate immune system to microbial infection is the ingestion, destruction and clearance of pathogens by activated phagocytic cells. Within their phagosomal compartments, macrophages generate an environment that is hostile to microbes via the production of reactive oxygen and nitrogen species, the elaboration of proteases, acidification of the phagosomal lumen and the export of iron.1,2 Unlike iron, copper (Cu) is not essential for many pathogens, although it is an essential trace element for mammals. In fact, Cu has been used for over a century as a microbiocidal agent for the elimination of human and plant pathogens.3 Moreover, Cu is critical for normal innate immune cell function and Cu deficiency in mammals renders the host susceptible to microbial pathogens. It is well documented that Cu levels rise dramatically in the serum of mammals during infection/inflammation and Cu accumulates at sites of inflammation, underscoring a specialized role for Cu in innate immunity.4,5 Consistent with a microbiocidal role for Cu in macrophages, phagosomal Cu concentrations rise dramatically in activated macrophages in parallel with increased expression of the Ctrl plasma membrane Cu importer and the Atp7A Cu pump on the phagosomal membrane.6, 7 Moreover, studies in both prokaryotic and fungal pathogens demonstrate the requirement for the Cu homeostasis machinery in resistance to macrophage killing.7, 8 Together, these studies provide compelling evidence for a critical role of Cu in microbiocidal activity of the host, which is countered by the Cu homeostasis machinery of invading bacterial and fungal pathogens.\nBiocidal properties of Cu have been documented since ancient Egyptians used it for water and wound sterilization in 2400 B.C.9 The Bordeaux mixture of copper sulfate and lime has been used since 1880 as a fungicide for grapes and other plants,3 while metallic Cu surfaces, which likely release Cu ions upon bacterial contact, may reduce contamination in hospitals.10-12 Metallic Cu, Cu salts, and Cu compounds continue to be used to control bacterial, fungal, and algal growth in agricultural and healthcare settings where microbes exist in the environment.3 The environment of pathogenic microbes, however, is ultimately the infected host, and while Cu and its complexes have been used for so long as environmental antimicrobials, no approach to date has been effective at using Cu to fight infection in a mammalian host. A challenge in this area is to discover ways to deliver or reallocate Cu selectively to the site of infection or inflammation.4 "}
{"text": "Conventionally, there is a medicine dispensing device for dispensing blister packages having a configuration in which a blister package is conveyed with the help of a grip unit and cut by a cutter mechanism to enable retrieval of the required quantity of the package sheet (see Patent document 1, for example).\nAs another medicine dispensing device for dispensing blister packages, there is one having a configuration that enables retrieval of a blister package loaded inside a medicine cassette by sucking it with a suction member (see Patent document 2, for example).\nHowever, in the medicine dispensing device disclosed in Patent Document 1, a grip unit is used for retrieving not only a fraction of the blister package but also one full blister package, and therefore, there is a problem of poor retrieving efficiency. Further, since the blister packages are stacked vertically, the space occupied in the height direction increases, and there is a problem that the quantity and type that can be accommodated is limited.\nIn the medicine dispensing device disclosed in Patent Document 2, it is not possible to dispense fraction of a blister package."}
{"text": "Plastic exterior products in a variety of colors and with a high quality sense of touch are increasingly popular for electronic parts, automobile parts and the like. Plastic exterior products can be formed of a resin composition including a thermoplastic resin and metal particles distributed therein to give a metal-like texture to the product.\nJapanese Patent Laid-Open Publication Nos. 2001-262003 and 2007-137963 discuss adding metal particles to a plastic resin, but the metal-like texture did not appear in an actual experiment. Also, there is a problem in that a flow mark or a weld line is formed during an injection molding process.\nJapanese Patent Laid-Open Publication No. 1995-196901 discusses adding a metal microplate having an average aspect ratio (average thickness/average particle diameter) of about 1/100 to about 1/8 obtained through a punching process to a resin to provide a metal-like texture. This process, however, is also limited by formation of a weld line."}
{"text": "Large quantities of glucose-containing syrups are manufactured by the enzymatic hydrolysis of corn starch. The first step of this process is usually accomplished by treating a starch-water mixture with an alpha-amylase enzyme at temperatures near the boiling point of water. A thermostable alpha-amylase is most efficient at these temperatures. For this reason, there is a great incentive to obtain a low-cost source of such a thermostable enzyme.\nThe alpha-amylases used commercially are produced by various microorganisms. It is known that these microorganisms contain genetic material which codes for the production of the enzymes by the organism. This genetic material is present in the form of deoxyribonucleic acid, hereafter referred to as DNA, within the cell.\nMost genetic material in a bacterium exists as giant DNA molecules which are present as the chromosome of the cell. However, a certain amount of the genetic material may also be present in the form of smaller closed circular molecules of DNA, known as plasmids.\nBy techniques referred to as genetic engineering, it is possible to transfer a portion of the DNA from one organism to another. Various workers have attempted to use these techniques to develop microorganisms which are superior alpha-amylase producers.\nOne technique that has been used is to remove the total DNA, i.e., chromosomal plus plasmid DNA, from a microorganism known to produce an amylase. Fragments of this DNA are then linked with the DNA of a bacteriophage. The bacteriophage containing this combined DNA is used to infect another microorganism whereby it transfers this genetic material onto the chromosomal DNA of the microorganism. The cells of the infected organism that receive DNA containing an amylase gene in an active form become amylase producers. These cells can be selected and grown for further use.\nAn additional genetic engineering technique that has been used with thermostable alpha-amylases involves extracting the total DNA from a microorganism known to produce such amylases. Some fragments of the DNA are transformed into a new host microorganism whereby they are incorporated into the chromosome of the host microorganism. Any cells into which the DNA containing active amylase-coding genetic material has been introduced are then selected and grown.\nIt has now been found that DNA coding for a thermostable alpha-amylase can exist in a naturally-occurring plasmid. Furthermore, it has been found that such naturally-occurring genetic material containing a thermostable alpha-amylase coding gene can be combined directly with another plasmid, known as a plasmid vector, to give a synthetic plasmid, hereafter called a chimeric plasmid. This chimeric plasmid can be inserted into and multiplied in a new host microorganism. This permits the development of microorganisms that are superior alpha-amylase producers which can be grown readily on a commercial scale. It is to the use of the newly-discovered, naturally-occurring plasmids containing alpha-amylase genes, the formation of chimeric plasmids containing their genetic material, and the development of new microorganisms containing these chimeric plasmids, that the present invention is directed."}
{"text": "The desirability of improving the efficiency of combustion in a vehicle's engine has long been recognized. For instance, Lyons and McKone is U.S. Pat. No. 2,086,775, and again in U.S. Pat. No. 2,151,432, disclose a method for improving combustion efficiency in an internal combustion engine by adding to the fuel what is described as \"relatively minute quantities\" of catalytic organometallic compounds. The Lyons and McKone patents, though, are directed solely to internal combustion engines and do not address the problem of NO.sub.x emissions from diesel engines.\nIn a unique application of catalytic technology described in International Publication No. WO 86/03482 and U.S. Pat. No. 4,892,562, Bowers and Sprague teach the preparation of diesel fuel containing fuel soluble platinum group metal compounds at levels of from 0.01 to 1.0 parts per million (ppm). The Bowers and Sprague results were corroborated and refined by work of Kelso, Epperly, and Hart, described in \"Effects of Platinum Fuel Additive on the Emissions and Efficiency of Diesel Engines,\" Society of Automotive Engineers (SAE) Paper No. 901,492, August 1990. Although the use of platinum group metal additives is effective, further nitrogen oxides reductions are still believed possible.\nMoreover, in \"Assessment of Diesel Particulate Control--Direct and Catalytic Oxidation,\" SAE Paper No. 81 0112, 1981, Murphy, Hillenbrand, Trayser, and Wasser have reported that the addition of catalyst metal to diesel fuel can improve the operation of a diesel trap. Among the catalysts disclosed is a platinum compound, albeit one containing chlorine, known to reduce catalyst effectiveness. In addition, the regeneration of a diesel trap by the use of a metallic additive which can include copper, nickel, cobalt, and, especially, iron, is discussed by Mueller, Wiedemann, Preuss and Schadlich in \"Diesel Particulate Filter System With Additive Supported Regeneration,\" ATZ Automobiltechnische Zeitschift 91 (1989).\nOther researchers have considered the use of water in oil emulsions for improving combustion efficiency in diesel engines. For instance, DenHerder, in U.S. Pat. No. 4,696,638, discusses such emulsions and indicates that the positive effects therefrom include \"cleaner exhaust.\" Although the disclosure of Denherder refers to emulsions containing up to about 40% water, Denherder is primarily directed to emulsions having only up to about 10% water in the form of droplets having a diameter of about 1 to about 10 microns.\nFurthermore, Wasser and Parry have reported in \"Diesel Engine NO.sub.x Control: Selective Catalytic Reduction and Methanol Emission,\" EPRI/EPA Joint Symposium on Stationary NO.sub.x Control, New Orleans, La., March, 1987 that NO.sub.x reductions of up to 80%, which are the levels desired for effective emission control, can be achieved in diesel engines using water and oil emulsions. They found, though, that emulsions of at least 60% water in oil are necessary to achieve such reductions. Unfortunately, such high water ratios can lead to increased emissions of carbon monoxide (CO) and unburned hydrocarbons. In addition, such high water levels can also cause problems in emulsion stability and create corrosion and storage volume concerns.\nRecently, Dainoff and Sprague, in U.S. patent application entitled \"Process for Reducing Nitrogen Oxides Emissions and Improving the Combustion Efficiency of a Turbine:, Ser. No. 07/691,556, filed Apr. 25, 1991, have discovered that water-in-fuel oil emulsions up to about 50% water by weight are useful for reducing nitrogen oxides and particulate emissions from combustion turbines, but they do not address the usefulness of such emulsions in vehicle diesel engines.\nThere has been a great deal of study of processes and agents for reducing nitrogen oxides emissions from stationery sources, such as boilers, etc. Preferred among such processes are those referred to as selective, non-catalytic reduction (SNCR) NO.sub.x reducing processes. Such processes involve the introduction of a NO.sub.x reducing treatment agent into a combustion effluent to achieve reductions of nitrogen oxides of up to about 50% or greater by a gas phase free radical mediated reaction process. SNCR processes generally utilize a nitrogen containing treatment agent such as ammonia, urea, or other chemicals and compositions which can break down in the effluent to the amidozine radical, which is believed to be the moiety which reacts with and reduces nitrogen oxides to molecular nitrogen (N.sub.2).\nSuch processes are disclosed by, for instance, Lyon, in U.S. Pat. No. 3,900,554, Arand, Muzio, and Sotter, in U.S. Pat. No. 4,208,386, and Arand, Muzio, and Teixeira, in U.S. Pat. No. 4,325,924. Although effective at reducing nitrogen oxides in effluents from open flame combustion, it has been believed that at the temperatures and pressures which exist in diesel engines, the use of nitrogenous nitrogen oxides reducing agents in such processes may actually create NO.sub.x, thus exacerbating the problem.\nAccordingly, a process and composition which is effective at substantially reducing the nitrogen oxides emissions from a diesel engine without the drawbacks of the prior art is extremely desirable."}
{"text": "P2 receptors have been generally categorized as either metabotropic nucleotide receptors or ionotropic receptors for extracellular nucleotides. Metabotropic nucleotide receptors (usually designated P2Y or P2Y.sub.n, where \"n\" is a subscript integer indicating subtype) are believed to differ from ionotropic receptors (usually designated P2X or P2X.sub.n) in that they are based on a different fundamental means of transmembrane signal transduction: P2Y receptors operate through a G protein-coupled system, while P2X receptors are ligand-gated ion channels. The ligand for these P2X receptors is ATP, and/or other natural nucleotides, for example, ADP, UTP, UDP, or synthetic nucleotides, for example 2-methylthioATP.\nAt least seven P2X receptors, and the cDNA sequences encoding them, have been identified to date. P2X.sub.1 cDNA was cloned from the smooth muscle of the rat vas deferens (Valera et al (1994) Nature 371:516-519) and P2X.sub.2 cDNA was cloned from PC12 cells (Brake et al. (1994) Nature 371:519-523). Five other P2X receptors have been found in cDNA libraries by virtue of their sequence similarity to P2X.sub.1 and P2X.sub.2 (P2X.sub.3 : Lewis et al. (1995) Nature 377:432-435, Chen et al. (1995) Nature 377:428-431; P2X.sub.4 : Buell et al. (1996) EMBO J. 15:55-62, Seguela et al. (1996) J. Neurosci. 16:448-455, Bo et al. (1995) FEBS Lett. 375:129-133, Soto et al. (1996) Proc. Natl. Acad. Sci. USA 93:3684-3688, Wang et al. (1996) Biochem. Biophys. Res. Commun.220:196-202; P2X.sub.5 : Collo et al. (1996) J. Neurosci. 16:2495-2507, Garcia-Guzman et al. (1996) FEBS Lett. 388:123-127; P2X.sub.6 : Collo et al. (1996), supra, Soto et al. (1996) Biochem. Biophys. Res. Commun. 223:456-460; P2X.sub.7 : Surprenant et al. (1996) Science 272:735-738). For a comparison of the amino acid sequences of rat P2X receptors see Buell et al. (1996) Eur. J. Neurosci. 8:2221-2228.\nNative P2X receptors form rapidly activated, nonselective cationic channels that are activated by ATP. Rat P2X.sub.1 and rat P2X.sub.2 have equal permeability to Na.sup.+ and K.sup.+ but significantly less to Cs.sup.+. The channels formed by the P2X receptors generally have high Ca.sup.2+ permeability (P.sub.Ca /P.sub.Na.congruent.4). The cloned rat P2X.sub.1, P2X.sub.2 and P2X.sub.4 receptors exhibit the same permeability for Ca.sup.2+ observed with native receptors. However, the mechanism by which P2X receptors form an ionic pore or bind ATP is not known.\nA variety of tissues and cell types, including epithelial, immune, muscle and neuronal, express at least one form of P2X receptor. The widespread distribution of P2X.sub.4 receptors in the rat central nervous system suggests a role for P2X.sub.4 -mediated events in the central nervous system. However, study of the role of individual P2X receptors is hampered by the lack of receptor subtype-specific agonists and antagonists. For example, one agonist useful for studying ATP-gated channels is .alpha.,.beta.-methylene-ATP (.alpha.,.beta.meATP). However, the P2X receptors display differential sensitivity to the agonist with P2X.sub.1 and P2X.sub.2 being .alpha.,.beta.meATP-sensitive and insensitive, respectively. Furthermore, binding of .alpha.,.beta.meATP to P2X receptors does not always result in channel opening. The predominant forms of P2X receptors in the rat brain, P2X.sub.4 and P2X.sub.6 receptors, cannot be blocked by suramin or PPADS. These two forms of the P2X receptor are also not activated by .alpha.,.beta.meATP and are, thus, intractable to study with currently available pharmacological tools.\nA therapeutic role for P2 receptors has been suggested, for example, for cystic fibrosis (Boucher et al. (1995) in: Belardinelli et al. (eds) Adenosine and Adenine Nucleotides: From Molecular Biology to Integrative Physiology (Kluwer Acad., Norwell Mass.) pp 525-532), diabetes (Loubatieres-Mariani et al. (1995) in: Belardinelli et al. (eds), supra, pp 337-345), immune and inflammatory diseases (Di Virgilio et al. (1995) in: Belardinelli et al. (eds), supra, pp 329-335), cancer (Rapaport (1993) Drug Dev. Res. 28:428-431), constipation and diarrhea (Milner et al. (1994) in: Kamm et al. (eds.) Constipation and Related Disorders: Pathophysiology and Management in Adults and Children (Wrightson Biomedical, Bristol) pp 41-49), behavioral disorders such as epilepsy, depression and aging-associated degenerative diseases (Williams (1993) Drug. Dev. Res. 28:438-444), contraception and sterility (Foresta et al. (1992) J. Biol. Chem. 257:19443-19447), and wound healing (Wang et al. (1990) Biochim. Biophys. Res. Commun. 166:251-258).\nAccordingly, there is a need in the art for specific agonists and antagonists for each P2X receptor subtype and, in particular, agents that will be effective in vivo, as well as for methods for identifying P2X receptor-specific agonist and antagonist compounds."}
{"text": "This invention relates to pneumatic tires and more particularly to a pneumatic tire having a wide reinforcing breaker or belt extending into an annular protuberance in the buttress region of the tire.\nThe tread reinforcement component commonly referred to as a belt or breaker is disposed between the tread and the carcass in the crown region of the tire. Typically, a breaker or belt is constructed of one or more plies of rubberized cord fabric. Each cord in the fabric is generally of a substantially non-extensible material such as metallic wires or strands, glass filaments, or rayon filaments. The cords within the breaker plies are oriented parallel to each other and substantially parallel to the plane of the beads and thus circumferentially of the tire, or at a small bias angle. In a single ply breaker construction, the small bias angle can be 0.degree.. Where the breaker is a multi-ply structure, however, the small bias angle is in the range of about 15.degree.-30.degree. with respect to the median equatorial plane of the tire, and may be oppositely disposed in successive plies.\nThe development of a reinforcing belt or breaker for use in pneumatic tires added a number of benefits to tires including considerable improvement to tire durability. On the other hand, the breaker also introduced some new problems of its own. The lateral edges of a breaker in a typical tire, due to their location between the tire tread and carcass at a position of continually changing stress, tend to separate from the carcass or themselves. Improper positioning of the breaker results in the uniformity problems of increased lateral force variation and conicity.\nAnother benefit of a reinforcing breaker is improved tread wear due to the greater stiffness imparted to the tire tread. Of course, this benefit is enhanced by utilizing a \"flat\" tread having no convex cross-sectional curvature so that the relative movement between the center portion of the tread and the axially outward portions is minimized thereby resulting in an even wear pattern. A serious drawback of a tread designed to have no convex cross-sectional curvature lies in the fact that a tire which is vulcanized with such a shape, when cured, may inflate to a concave tread cross-section rather than a flat cross-section. Elimination of this problem can be accomplished by constructing and vulcanizing the tire such that the tread has at least a small convex cross-sectional curvature, that is, a large crown radius ratio, when the cured tire is inflated to its rated inflation pressure.\nIn the detailed description of the invention following hereinafter, the cross-sectional curvature of the tread and breaker of the tire are discussed in terms of \"crown radius ratio\" and \"crown radius.\" The term \"crown radius\" as used herein either with respect to the tread or breaker, means the radius of the arc of a circle which best approximates the arc of curvature of the tread surface or the breaker as the case may be of a tire. The crown radius is commonly determined by the formula ##EQU1## where \"leg set\" is the axial distance measured along a line tangent to the tread or breaker at its midpoint between the shoulders of the tread or edge of the breaker, or predetermined points axially equidistance inward of the shoulders or breaker edges. Shoulder drop is the perpendicular distance from the aforementioned tangent line to the tread shoulders or said predetermined points. In the case of a breaker, shoulder drop may be considered as the perpendicular distance between the line tangent to the breaker at its midpoint and the edges of the breaker. The \"crown radius ratio\" is the ratio of the crown radius of either the tread or breaker to the nominal cross-sectional width of the tire and as used herein is expressed in terms of percent.\nEither a belted radial or a belted bias tire construction is suitable for use in the present invention. The radial type tire is, however, the preferred embodiment. A radial tire is inclusive of various tire constructions which typically comprise a carcass or body having one or more reinforcement plies of cord fabric extending from bead to bead wherein the cords in each ply are substantially radial in orientation, that is, the cords are oriented substantially normal to the beads and the crown centerline of the tire. In a single ply radial tire construction, the carcass cords normally have a 90.degree. bias angle, that is, in the unshaped carcass they extend perpendicular to the planes of the beads. In a two-ply radial tire construction, the cords in each carcass ply are parallel to the cords in the other carcass ply. However, the cords in each ply may be oriented at oppositely disposed angles of 70.degree. or more with respect to the median equatorial plane of the tire and thus the angle between cords in different plies is between 0.degree. and 40.degree.. In more than two ply radial ply tire constructions, similar cord arrangements in successive carcass plies is usually employed.\nWith regard to the prior art, U.S. Pat. No. 3,450,182 discloses a tire having a flat tread cross-section, wide protuberances extending beyond the width of the tire sidewalls, a tread surface in normal engagement with the ground between the lateral edges of the protuberances and a breaker extending laterally beneath the tread to positions adjacent the edges of the protuberances. This patent does not, however, disclose a breaker having its lateral edges removed from the stress area of contact by the tread with the ground or a convex cross-sectional tread curvature facilitating the manufacture of a tire. Moreover, the extreme lateral extension of the protuberances beyond the tire sidewall will result in a high degree of tread shoulder wear and breaker separation due to centrifugal force moving the protuberances radially outward during tire rotation.\nIn U.S. Pat. No. 2,477,754, an aircraft tire is shown in which a tread and breaker have a substantial convexly curved cross-section when the tire is inflated and the breaker has edges positioned within protuberances extending from the side of the tire. When this tire is loaded (but not impacted), the tread-ground contact area is in the crown area of the tread only and consequently the tread wear and traction of the tire will be relatively poor. Moreover, when the tire is subject to the normal type of impaction and very heavy loading which aircraft tires receive, the lateral edges of the breaker rotate through the tread-ground contact area to cause edge separation of the breaker. A further defect in the tire of U.S. Pat. No. 2,477,754 is that the circumferential breaker cords along the lateral sides of the breaker will carry the entire breaker load under deflection so that these breaker cords will break after a small amount of such use."}
{"text": "Collapsible tube squeezing devices for dispensing the contents of tubes containing toothpaste, hair creme, hand lotion, shoe polishing creme, glue, and the like are well known in the art. A review of the prior art devices shows that while such devices exist they have not received substantial user acceptance due to their expense, size, lack of universality, difficulties with operation and the necessity to use both of the operator's hands in emptying the contents of collapsible tubes.\nThe advent of collapsible tubes formed of plastic, rather than the well known metallic type collapsible tubes, has also presented additional problems which have rendered many pre-plastic tube squeezing devices unsuitable for their intended purposes particularly in regard to removal of the contents of the crown portion of such tubes. Crown portions of plastic collapsible tubes have a tendency to be substantially more rigid than metallic tubes which factor coupled with the resiliency of plastics often employed in the manufacture of plastic collapsible tubes necessitates that special measures be undertaken in order to completely empty such tubes."}
{"text": "The subject matter disclosed herein relates to a rotary device and, more particularly, to a rotary device with no-back and torque-limiting capabilities.\nAn aircraft utilizes actuation systems for control of aircraft flight surfaces. Such actuation systems transmit torque to actuators through torque transmission devices including torque tubes, gear boxes and bearing supports. If a disconnection failure occurs to a torque transmission device, a no-back will prevent a loss of position control for a given flight surface by grounding the resultant torque generated by airloads on the flight control surface to a structural ground and thereby lock the surface in a fixed position.\nWith such flight control actuation systems, it is often necessary to mitigate an overload condition generated by a structural jam. One method to accomplish overload mitigation is to have a torque limiting device within the actuation system such that the actuator output torque is limited by a mechanism within an actuator by transferring torque to structural ground if the limit is exceeded."}
{"text": "1. Field of the Invention\nThe invention relates to a digital switching stage having a differential amplifier with a first and a second differential amplifier branch; a first resistor which is connected in the first differential amplifier branch has a first terminal that is a terminal for a first supply potential and has a second terminal; and a first bipolar transistor which is connected in an emitter follower circuit with respect to the second terminal of the first resistor has an emitter which is connected to an output terminal.\nSuch a switching stage is shown in the literature, such as in an article by George R. Watson, entitled \"Advances in Bipolar VLSI\", in Proceedings of the IEEE, Vol. 78, No. 11, pp. 1706-1719, November 1990. An output of the switching stage is essentially capacitively loaded by a connecting line connected to it. At a positively oriented edge of the output signal, the current required to reverse the charge of the parasitic capacitance which is operative on the output side flows through the emitter follower transistor. That is, the charge reversing current is furnished at low impedance by the first supply potential source. At a negatively oriented output signal edge, the charge reversing current flows either through one additional current source or through the input of a further switching stage connected to the output side. The current is furnished at high impedance from a further supply potential source.\nIn order to achieve a high operating speed, the switching stages must have adequate driver capability and a high permanent current therefore flows through the additional current source or the following switching stage. The power loss is correspondingly high.\n2. Summary of the Invention\nIt is accordingly an object of the invention to provide a digital switching stage with a current switch, which overcomes the hereinafore-mentioned disadvantages of the hereto-fore-known devices of this general type and which has a high processing speed and a low power loss.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, a digital switching stage, comprising a differential amplifier having a first and a second differential amplifier branch; a terminal for a first supply potential, a terminal for a second supply potential, and an output terminal; a first resistor being connected in the first differential amplifier branch and having a first terminal and a second terminal, the first terminal being the terminal for the first supply potential; a first bipolar transistor being connected in an emitter follower circuit with respect to the second terminal of the first resistor and having an emitter being connected to the output terminal; a second bipolar transistor having a base and having a collector-to-emitter path being connected between the output terminal and the terminal for the second supply potential; and a second resistor being connected in the second differential amplifier branch and having a first terminal and a second terminal, the first terminal of the second resistor being connected to the output terminal, and the second terminal of the second resistor being connected to the base of the second bipolar transistor.\nIn the switching stage according to the invention, for both the leading and the trailing output signal edges, the charge reversing current of the parasitic capacitance that is operative on the output side is furnished at low impedance from the supply potential forces. The input and output signal levels are symmetrical to the supply potentials. In comparison with the conventional switching stage, the supply voltage is lower, so that the lower loss is further reduced thereby.\nIn accordance with another feature of the invention, there is provided a current source transistor having a collector-to-emitter path, the differential amplifier including a first and at least one second switching transistor having mutually coupled emitters being connected through the collector-to-emitter path of the current source transistor to the terminal for the second supply potential.\nIn accordance with a further feature of the invention, there is provided a voltage divider connected between the terminals for the supply potentials for supplying the base of the first switching transistor with a reference potential.\nIn accordance with an added feature of the invention, the voltage divider includes a first diode having an anode connected to the terminal for the first supply potential and having a cathode; a second diode having a cathode connected to the terminal for the second supply potential and having an anode; a first resistor connected between the cathode of the first diode and the base of the first switching transistor; and a second resistor connected between the base of the first switching transistor and the anode of the second diode.\nIn accordance with an additional feature of the invention, the base of the current source transistor is connected to the anode of the second diode.\nIn accordance with yet another feature of the invention, the first and second resistors of the voltage divider have the same given resistance.\nIn accordance with yet a further feature of the invention, the first and second resistors of the differential amplifier each have the same resistance being twice the given resistance.\nIn accordance with a concomitant feature of the invention, the resistance of the first resistor of the differential amplifier is greater than the resistance of the second resistor of the differential amplifier, and the resistances of the first and second resistors of the differential amplifier are each greater than twice the given resistance.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in a digital switching stage with a current switch, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims."}
{"text": "Cancer generally refers to one of a group of more than 100 diseases caused by the uncontrolled growth and spread of abnormal cells that can take the form of solid tumors, lymphomas, and non-solid cancers such as leukemia. Unlike normal cells, which reproduce until maturation is attained and then only as necessary for replacement, cancer cells grow and divide endlessly, crowding out nearby cells and eventually spreading to other parts of the body, unless their progression is stopped. Once cancer cells metastasize by leaving a tumor, they will travel through the bloodstream or lymphatic system to other parts of the body, where the cells begin multiplying and developing into new tumors. This sort of tumor progression makes cancer dangerously fatal. Although there have been great improvements in diagnosis, general patient care, surgical techniques, and local and systemic adjuvant therapies, most deaths from cancer are still due to metastases and other cancers that are resistant to conventional therapies including radiation and chemotherapy.\nRadiation therapy is typically only effective for cancer treatment at early and middle stages of cancer, when cancer is localized, and not effective for late stage disease with metastasis. Chemotherapy can be effective at all stages of the disease, but there can be severe side effects, e.g. vomiting, low white blood cells, loss of hair, loss of weight and other toxic effects, to radiation therapy and chemotherapy. Because of such severe side effects, many cancer patients do not or cannot successfully complete a chemotherapy treatment regimen. The side effects of radiation and anticancer drugs can be viewed as resulting from poor target specificity. Anticancer drugs, typically administered intravenously or more rarely orally, circulate through most normal tissues of patients as well as the target tumors. If the drug is toxic to a normal cell, then this circulation will result in the death of normal cells, leading to side effects, and the more toxic the drug to normal cells, the more serious the side effects. Due to these and other problems, some highly cytotoxic chemotherapeutic agents, agents with nanomolar or sub-nanomolar IC50 values against cancer cells, have not been successfully developed into approved drugs.\nProdrugs have been investigated as a means to lower the unwanted toxicity or some other negative attribute of a drug without loss of efficacy. A prodrug is a drug that has been chemically modified to render it inactive but that, subsequent to administration, is metabolized or otherwise converted to the active form of the drug in the body. For example, in an effort to improve drug targeting, prodrugs have been developed that are activated under hypoxic conditions. Hypoxia creates a bioreductive environment, and certain anti-cancer agents have been converted into prodrugs that can be activated in such environments. See the reviews by Naylor et al., May 2001, Mini. Rev. Med. 1(1):17-29, and Denny, 2001, Eur. J. Med Chem. 36: 577-595. “Hypoxia” is a condition of low oxygen levels; most solid tumors larger than about 1 mm in diameter contain hypoxic regions (see the references Coleman, 1988, J. Nat. Canc. Inst. 80: 310; and Vaupel et al., Cancer Res. 49: 6449).\nAs a tumor grows, it requires a blood supply and thus the growth of new vasculature. The new vasculature that supports tumor growth is often highly unordered, leaving significant portions of the tumor under-vascularized and subject to intermittent vascular blockage. The vascular architecture of the tumor can contribute significantly to the cancer's ability to survive drug therapy in at least two different ways. First, if the drug must reach the cancer through the bloodstream, then not as much drug will reach the under-vascularized, hypoxic areas of the tumor. Second, to the extent the drug requires oxygen to be effective, then the drug will be less effective in the hypoxic regions of the tumor.\nConversely, however, the hypoxic environment is conducive to reductive events that can be used to generate reduced derivatives of a variety of chemical groups (see the reference Workman et al., 1993, Cancer and Metast. Rev. 12: 73-82), and bioreductive prodrug compounds have been developed to exploit such environments. These prodrugs include the antibiotics Mitomycin C (MMC) and Porfiromycin (POR), N-oxides such as Tirapazamine (TRZ; see the reference Zeeman et al., 1986, Inst. J. Radiot. Oncol. Biol. Phys. 12: 1239), quinones such as the indoloquinone E09 (see the reference Bailey et al., 1992, Int. J. Radiot. Oncol. Biol. Phys. 22: 649), cyclopropamitosenes (EP-A-0868137), and a tertiary amine-N-oxide analogue of Mitoxantrone (AQ4N) that is activated by cytochrome P450 3A4 (see the references Patterson, 1993, Cancer Metast. Rev. 12: 119; and Patterson, 1994, Biochem. Pharm. Oncol. Res. 6: 533).\nOther bioreductively activated prodrug compounds include the nitroimidazole derivatives that have been reported to be useful in cancer radiotherapy as radio-sensitizing agents (see the patent publications EP312858 and WO91/11440) and potentiatiors of chemotherapeutic agents (see U.S. Pat. No. 4,921,963). Nitroimidazole has also been conjugated to the anti-cancer agent PARP 5-bromoisoquinolinone (see the reference Parveen et al., July 1999, Bioorg. Med. Chem. Lett., 9:2031-36). The nitroimidazole moiety itself is, however, somewhat cytotoxic to normal cells, because it undergoes redox cycling and generates superoxides under oxygenated conditions.\nBioreductively activated prodrug compounds that include a nitroimidazole linked to a variety of anti-neoplastic agents has been described in A 2-NITROIMIDAZOLE CARBAMATE PRODRUG OF 5-AMINO-1-(CHLOROMETHYL)-3-[5,6,7-TRIMETHOXYINDOL-2-YL)CARBONYL]-1,2-DIHYDRO-3H-BENZ[E]INDOLE (AMINO-SECO-CBI-TMI) FOR USE WITH ADEPT AND GDEPT, M. P. Hay et al., Bioorganic & Medicinal Chemistry Letters 9 (1999) 2237-2242, and PCT publication WO 00/64864. In all of the prodrugs described in these documents, the nitroimidazole is directly linked to a carbamate linker and the anti-neoplastic agent is protected at a nitrogen via a carbamate or at a carbon via an ester linkage.\nThus, there remains a need to provide drugs to treat cancer. Such drugs would be especially beneficial if they targeted cancer cells more effectively than current drugs and had fewer, less serious side effects. The present compounds and methods help meet this need."}
{"text": "Fiber optics communication systems require compact light emitting sources capable of generating single-mode, tunable, narrow linewidth radiation in the 1.3-1.56 .mu.m wavelength range. Some of the existing semiconductor lasers, for example, InGaAsP DFB lasers can meet requirements for high power and proper wavelength, but a high dynamic single mode yield is difficult to achieve. Conventional index coupled CFB lasers employing an index corrugation have an inherent problem in existence of two longitudinal modes with an equal threshold gain which results in poor single mode operation as shown, for example, in the article by H. Kogelnik and C. V. Shank \"Coupled-Wave theory of distributed feedback lasers\", J. Appl. Phys., vol. 43, no. 5, pp. 2327-2335, 1972. In complex coupled DFB lasers, a periodic optical gain or loss modulation in the presence or absence of conventional index corrugation along the laser cavity effectively breaks the mode degeneracy between the two Bragg modes around the stop band of the DFB lasers, and thus avoids a serious and inherent problem for conventional index coupled DFB lasers, as shown in publications by Y. Luo, Y. Nakano, K. Tada et al. \"Purely gain-coupled distributed feedback semiconductor 35 lasers\", Appl. Phys. Lett., vol. 56, pp. 1620-1622, 1990 and G. P. Li, T. Makino, R. Moore et al. \"Partly gain-coupled 1.55 .mu.m strained layer multi-quantum well DFB lasers\", IEEE J. Quantum Electronics, vol. QE-29, pp. 1736-1742, 1993. With introduction of even a small amount of gain or loss coupling, the dynamic single mode yield of complex coupled DFB lasers increases drastically, whether or not there is index coupling. It effectively provides lasing predominantly on a preferred and fixed Bragg mode among the two originally degenerate ones around a stop band, regardless of random distribution of unknown laser facet phases. In in-phase gain coupled DFB lasers, the higher index region within a grating period has a higher optical modal gain, resulting in lasing mainly of the right Bragg mode with the lasing wavelength longer than the Bragg wavelength. Both theory and experiments have confirmed this conclusion. For anti-phase loss coupled DFB lasers, the higher index region within a grating period experiences higher optical loss, resulting primarily in lasing of the left Bragg mode with the lasing wavelength shorter than the Bragg wavelength.\nPreviously known gain coupled or loss coupled DFB lasers utilize only one type of complex coupling which results in a high dynamic single mode yield for one of predetermined and dominant Bragg mode only. The random distribution of cleaved or HR-coated laser facets is usually the major reason why the opposite Bragg mode across the stop band sometimes actually lases as the dominant mode. Nevertheless, with an increase of complex coupling strength, the single mode yield increases quickly, and an excellent dynamic single mode behavior with a high side mode suppression ratio (SMSR) can be obtained, as reported in the publication by B. Borchert, K. David, B. Stegmuller et al. \"1.55 .mu.m gain-coupled quantum-well distributed feedback lasers with high-single-mode yield and narrow linewidth\", IEEE Photon. Technol. Lett., vol. 3, no. 11, pp. 955-957, 1991.\nComplex coupled DFB lasers have been studied intensively by different scientific groups. It was shown that they constitute a good candidate for practical applications in fiber optics communication systems. Some of the achievements are reported in the following publications: W. T. Tsang, F. S. Choa, M. C. Wu et al. \"Semiconductor distributed feedback lasers with quantum well superlattice gratings for index or gain-coupled feedback\", Appl. Phys. Lett., vol. 60, pp. 2580-2582, 1992 and J. Hong, K. Leong, T Makino et al. \"Impact of random facet phase on modal properties of partly gain-coupled DFB lasers\", J. Selected Topics on Quantum Electronics, vol. 3, no. 2, pp. 555-568, 1997.\nNevertheless, a rapid advance in high speed and large capacity dense wavelength division multiplexing (DWDM) fiber optics systems has been constantly demanding semiconductor lasers having not only single mode properties but also a wide continuous tuning range for practical and cost effective applications. A number of techniques have been used in attempt to obtain a widely tunable DFB laser with dynamic single mode performance comparable with standard DFB laser generating at fixed wavelength. First, conventional index coupled tunable lasers followed by complex coupled lasers were tried. The latter appeared to be better in comparison with index coupled lasers because of maintaining good single mode properties within a wide wavelength tuning range. However, similar to other single contact DFB lasers, a wavelength range of complex coupled DFB lasers is limited either by the injection current or base temperature. A multi-sectional laser with separate electrodes for independent current injection was introduced next, providing an effective mode selection and phase control. Some of the experimental implementations of the idea are described in the following publications: Y. Yoshikuni, K. Oe, G. Motosugi, T. Matsuoka \"Broad Wavelength tuning under single-mode oscillation with a multi-electrode distributed feedback laser\", Electronic Lett., vol. 22, no. 22, pp. 1153-1154; I. Kim, R. C. Alferness, U. Koren et al. \"Broadly tunable vertical-coupler filtered tensile-strained InGaAs/InGaAsP multiple quantum well laser\", Appl. Phys. Lett., vol. 64, pp. 2764-2766, 1994; M. Oberg, S. Nilsson, K. Streubel et al. \"74 nm wavelength tuning range of an InGaAsP/In vertical grating-assisted codirectional coupler laser with rear sampled grating reflector\", IEEE Photonics Technol. Lett., vol. 5, pp. 735-737, 1993; H. Ishii, Y. Tohmori, Y. Yoshikuni et al. \"Multiple-phase-shift super structure grating DBR lasers for broad wavelength tuning\", IEEE Photonics Technol. Lett., vol. 5, pp. 613-615, 1993; M. Okai, I. F. Lealman, L. J. Rivers et al. \"In-line Fabry-Perot optical waveguide filter with quasi-chirped gratings\", Electron. Lett., vol. 32, pp. 108-109, 1996. When conventional index coupled DFB lasers are used in conjunction with multi-electrodes to achieve wider wavelength tuning range, current injection into two different electrodes of DFB lasers has to be well controlled and the current ratio among the electrodes needs to be precisely maintained. Due to the inherent Bragg mode degeneracy in conventional index coupled DFB lasers, each DFB section has its own stop band. The two Bragg modes around the related stop band compete equally for lasing as the dominant mode. Which one of the Bragg modes will tend to lase depends largely on the unknown grating facet phase and the injection currents in the two different sections. When two sections of the laser are biased independently, two pairs of Bragg modes, four in total, tend to compete against each other. Normally it results in a low SMSR or fast degradation of SMSR during the wavelength tuning within a certain current injection range. To overcome these problems, a complicated and tedious injection current control is normally required. As a result, a high yield device with a high SMSR, wide controlled tuning range and long term stable operation is difficult to obtain. Moreover, from device to device, the random facet phase usually causes a random distribution of lasing mode between the two pairs of degenerated Bragg modes, and the final single mode yield with a wide wavelength tuning range is low. Even if a preselected SMSR is initially high, it can degrade rapidly with a small injection current variation within a certain range due to the sensitivity of index coupled DFB lasers to initial phase variations.\nAmong other multi-sectional lasers, complex coupled lasers appear to be the best candidates for a wide range tuning. Complex coupled DFB lasers are preferably used as building blocks for tunable lasers to ensure a high SMSR and provide high single mode yield and a small degradation of SMSR over injection current variation. Nevertheless, the tuning range of these lasers is still limited which results in a continued demand for further increasing a tuning range of DFB lasers while maintaining other essential characteristics of the lasers, and development of effective methods of their operation."}
{"text": "1. Field of the Invention\nThe present invention relates to a laser light amount control method of a semiconductor laser scanner used for a writing head of a laser beam printer and a digital copying machine, particularly relates to a laser light amount control method in a tri-level developing system of forming a normal latent image and a reverse latent image by one exposure to execute two-color developing.\n2. Description of the Related Art\nIn a conventional laser printer, a method of developing one color of toner by one laser scanning beam is used. Recently, for one of methods for the color laser printer, a tri-level developing method is developed, in which a method is used for forming by one laser scanning beam a normal developing latent image, a reverse developing latent image and an intermediate potential latent image in which neither a normal developing nor a reverse developing is executed, thereby developing two colors of toner at a time.\nIn an electrophotographic process, when an electrostatic latent image is formed on a photosensitive drum, an electric field polarized opposite to the polarity of the latent image is generated in the circumference of the latent image in addition to an electric field for highlighting developing generated at the end of the latent image. In the conventional one-color developing, when an image is formed, field reversing does not matter. However, in the tri-level process in which positive and negative electrostatic latent images are formed on the photosensitive drum and the developing of two colors of toner such as red toner and black toner polarized opposite to the polarity of the respective latent images is executed, special developing (hereinafter called fringe developing) occurs that red toner adheres around a black image and black toner adheres around a red image because of field reversing. Therefore, a problem arises that the quality of a two-color image is remarkably deteriorated.\nThe fringe phenomenon remarkably emerges in a solid image. Though the situation of the emergence is different at the top end and at the back end in the rotational direction of the photosensitive drum, the fringe phenomenon particularly arises at the back end. Heretofore, efforts to reduce a fringe at the back end have been made.\nFor a method of reducing the fringe, there is a method of reducing the fringe by using various exposure controls for trade-off relations among the fringe, the adhesion of a carrier and a printing density, which is a theme of tri-level developing. Specifically, there is a method of controlling an electric potential field in steps around an image by controlling the exposure amount to improve the surface of the photosensitive drum.\nThe fringe has been reduced by controlling the exposure amount under the aforementioned background. In case of the reduction of the fringe, as described above, the exposure amount control has been focused at the back end of the solid image in which the fringe phenomenon is remarkable.\nHowever, as the experiments are made, it has proved that another correction is also necessary for a new fringe phenomenon arising between close vertical lines because the fringe phenomenon is remarkable not only at the back end of the solid image but between the close vertical lines.\nMoreover, in case that a reverse image and a normal image are adjacent in the horizontal scanning direction, the correction control generates an unnecessary image turbulence in a portion to which the correction exposure is applied under a solid part under a part between the adjacent images."}
{"text": "1. Field\nExemplary embodiments of the present disclosure relate to an apparatus and method for generating a medical image to realistically display an object or inner tissues of the object.\n2. Description of the Related Art\nResearch into medical imaging apparatuses is actively being conducted to go along with an increasing interest in improving the health of people. Examples of medical imaging apparatuses may include an X-ray imaging apparatus, an X-ray fluoroscopy apparatus, a computerized tomography (CT) scanner, a magnetic resonance imaging (MRI) apparatus, a positron emission tomography (PET) apparatus, and an ultrasonic diagnostic apparatus.\nThese medical imaging apparatuses display a 2D or 3D medical image of an object. A 2D medical image refers to a sectional image of inner tissues of an object. A 3D medical image refers to an image acquired through volume rendering of 3D volume data generated based on a plurality of sectional images.\nThe 2D and 3D medical images may be grayscale images or color images. Grayscale images are limited in the ability to provide a realistic image of an object, and thus, color images are currently the norm. Such color images may be acquired by mapping a grayscale image with colors that are similar to those of inner tissues of an object.\nHowever, colors used in mapping are arbitrarily selected and thus, there remain limitations in providing a more realistic color image."}
{"text": "BNP and proBNP are reliable markers of heart failure (HF) widely used in clinical practice. Several types of sandwich immunoassays (conventional assays) utilizing two mono- or polyclonal antibodies, specific to different epitopes of BNP or BNP-fragment of proBNP molecule are described in literature.\nBNP molecule is known as an extremely unstable molecule rapidly losing its immunological activity in water solutions. This loss of activity is usually associated with proteolytic degradation of the peptide. Sandwich immunoassays commonly used for qualitative or quantitative antigen immunodetection utilize two or more antibodies specific to two or more different epitopes. The longer is the distance between the epitopes, the higher is the probability that sites of proteolysis would be located between the epitopes of the antibodies, thus increasing the sensitivity of the assay to proteolytic degradation of the antigen. And vice versa, the closer are the epitopes to each other, the smaller is the probability of the proteolytic cleavage of the molecule between the epitopes.\nImmunoassay methods for very small molecules have been described, including the application of so called anti-metatype antibodies. Such methods are disclosed, e.g. for detecting digoxin (Self et al., 1994, Clin. Chem. 40:2035-2041), and angiotensin II (Towbin et al., 1995, J. Immunol. Meth. 181:167-176).\nHowever, it is not an easy task to apply this type of method to different analytes, since very specific monoclonal antibodies are required in such a method."}
{"text": "The present invention relates to a fluid control valve device suitably used as a rotary fluid channel switching valve in a power steering apparatus of a hydraulic type or the like for reducing a steering wheel operating force of, e.g., an automobile and, more particularly, to an improvement in a rotary fluid control valve device for controlling distribution of a fluid pressure by relative rotational displacement between a rotor and a sleeve.\nA fluid control valve device of this type is used as a fluid channel switching valve in a power steering apparatus or the like to operate a power cylinder as a movable unit in accordance with a steering wheel operation, thereby generating an auxiliary steering force during a steering operation. Various types of such a fluid control valve device have been conventionally proposed as disclosed in, e.g., Japanese Patent Laid-Open No. 57-178971 and Japanese Utility Model Publication No. 63-30613. That is, in this fluid channel switching valve, a rotor formed integrally with an input shaft (stab shaft) of a steering wheel and a sleeve formed integrally with its output shaft (pinion shaft) are assembled to be rotationally displaced relative to each other and arranged in a valve housing. A plurality of channel grooves formed in the outer circumferential surface of the rotor and the inner circumferential surface of the sleeve in their circumferential direction are connected to channels communicating with an oil pump as a fluid pressure generating source, an oil tank, and right and left cylinder chambers constituting a power cylinder, thereby easily and properly switching fluid channels of a fluid pressure circuit (hydraulic circuit). In addition, the arrangement of the device is simple.\nOne of problems posed when the above rotary fluid control valve device is used as a fluid channel switching valve in a power steering apparatus is a noise problem. That is, in the hydraulic circuit extending from the pump to the right and left cylinder chambers or to the tank through the pressurized oil channels in the fluid channel switching valve, an abrupt increase or decrease in a fluid channel sectional area forms a turbulence, a vortex, or the like in a flow of a pressurized oil in the pressurized oil channels constituted by the channel grooves or land portions formed in the rotor and the sleeve of the control valve and channel holes formed to open in the grooves or the land portions, thereby generating a fluid sound (so-called shoo sound). More specifically, when a wide or narrow portion, a bent portion, a branch portion, or the like is present in the fluid channel through which the pressurized oil flows, since a pressure variation is caused to form a turbulence or a vortex, generation of the fluid sound is more or less inevitable. However, this fluid sound is desired to be minimized in the fluid channel switching valve of this type.\nEspecially in the above rotary fluid channel switching valve for use in a power steering apparatus, the channel grooves and the land portions formed in the sleeve inner circumferential surface and the rotor outer circumferential surface are selectively combined by the relative rotational displacement therebetween, thereby connecting or disconnecting the fluid channels to perform switching between the fluid channels. Throttle portions formed between side edge portions of the land portions for controlling connection/disconnection of the fluid channels abruptly change the fluid channel sectional area because the side edge portions are formed as sharp corner portions. Therefore, a turbulence or a vortex is inevitably formed in the flow of a pressurized oil flowing from an input port of the pump into the channel grooves. For this reason, Japanese Utility Model Publication No. 63-38138, for example, proposes an arrangement in which the throttle portion defined between the two side edges of the channel grooves communicating with the input port or the two side edges of the land portions is formed to have a structure which does not cause much change in fluid channel sectional area by forming the shape of the side edge portion into an inclined surface shape, thereby solving the noise problem. However, this arrangement is still unsatisfactory to prevent the fluid sound caused by the flow of a pressurized oil. Therefore, a demand has arisen for a countermeasure capable of solving the above noise problem with a comparatively simple and inexpensive arrangement.\nIn addition, this noise problem similarly arises in channel grooves or a fluid discharge hole portion communicating with a return port to the tank in the above rotary fluid channel switching valve. Therefore, this point must also be considered."}
{"text": "1. Field of the Invention\nThe present invention relates to an information processing apparatus, a control method therefor, and a computer-readable storage medium storing a program for implementing the control method.\n2. Description of the Related Art\nIn recent years, as image forming apparatuses have become increasingly sophisticated, it has become common to store user data in a nonvolatile device such as an HDD (hard disk drive) or flash memory and read and use the user data when using a certain function.\nIn this case, when a nonvolatile device in which user data is stored fails, it becomes impossible to read the user data.\nAs a countermeasure against data loss caused by such a nonvolatile device failure, there is a mirroring function (RAID1) of constantly backing up data using two HDDs and providing notification when the HDDs fail.\nThere is also a failure predicting function of, before an HDD fails, using S.M.A.R.T. (Self-Monitoring, Analysis and Reporting Technology) information on the HDD to provide notification of a failure prediction indicating that a failure of the HDD is anticipated, thus urging a user to back up data.\nIn a case where, out of the mirroring function and the failure predicting function described above, only the mirroring function is enabled, data is constantly backed up, and hence a service engineer has only to replace the HDD after the HDD fails.\nUsing an HDD to the limit as described above has a merit of, for a user, eliminating downtime of an image forming apparatus caused by occurrence of an error, and a merit of, for the service engineer, reducing maintenance costs.\nOn the other hand, in a case where only the failure predicting function is enabled, upon receiving notification of a failure prediction, a service engineer needs to replace an HDD after backing up user data so as to ensure the user data. This failure prediction, however, is merely a prediction, and hence the HDD may fail immediately after notification or may be usable for several years after notification.\nIn a case where both the mirroring function and the failure predicting function are enabled, when notification of a failure prediction using the failure predicting function is provided, notification of an error is provided to a user although an HDD can still be used.\nAs a result, it is necessary to replace the HDD of an image forming apparatus for which notification of the error has been provided, resulting in undesired downtime and cost of maintenance by a service engineer.\nMoreover, when both the mirroring function and the failure predicting function are enabled, this means that an image forming apparatus is equipped with two HDDs. In this case, when notification of a failure prediction with respect to both the two HDDs is provided, it is necessary to notify a user of the failure prediction although the two HDDs are unlikely to fail at the same time. The user, however, does not know what to do even when he/she is notified of the failure prediction.\nIn relation to the above described technique, there has been disclosed a technique that S.M.A.R.T. information is stored in a server in case that the S.M.A.R.T. information is corrupted, and in the server, the S.M.A.R.T. information is compared with a previous log to thus reduce the load for analyzing the cause of a failure (see, for example, Japanese Laid Open Patent Publication (Kokai) No. 2001-312375).\nAlso, there has been disclosed a technique that a failure is predicted based on, for example, the number of accesses to an HDD or the like (see, for example, Japanese Laid Open Patent Publication (Kokai) No. 2006-256251).\nFor example, according to the technique described in Japanese Laid Open Patent Publication (Kokai) No. 2001-312375, a server can be notified of a status of an image forming apparatus, and also, values obtained using the failure predicting function can be backed up to the server.\nAccording to the technique described in Japanese Laid Open Patent Publication (Kokai) No. 2006-256251, a failure can be predicted based on the number of accesses. These conventional arts can improve the accuracy of the failure predicting function.\nHowever, when both the mirroring function and the failure predicting function are enabled, problems described hereafter arise. First, unnecessary notification of a failure prediction is provided to a user. This means that, for example, an image forming apparatus provides notification of a failure prediction to a user interface or the like.\nNext, unnecessary notification of a failure prediction is provided to a service engineer. This means that, for example, unnecessary notification of an error is provided to a server that manages an image forming apparatus, and a service engineer is called.\nMoreover, notification of a failure prediction is provided which causes a service engineer to falsely recognize that a trouble is serious and urgent. This means that, for example, notification of an error is provided to a server that manages an image forming apparatus, and the service engineer is called although a visit of the service engineer is unnecessary.\nFurther, even if a display indicative of a failure prediction notification is produced on a user interface of an image forming apparatus or the like, a user does not know what to do.\nAs described above, the conventional techniques have the problem of requiring unnecessary operational costs for calling a service engineer, replacing an HDD, and so on."}
{"text": "For hundreds of years, the use of certain non-toxic agents such as herbal compounds has been widely accepted for a variety of physiological conditions, especially in the Orient. Panax ginseng C.A. Meyer is the best known traditional Chinese medicine. The important pharmacological activities of ginseng extracts, alone or in combination with other drugs, include alleviation of renal impairment, inhibition of carcinogenesis and prevention of stress. There are also a number of reports on the influence of ginseng on the immunological responsiveness of the individual. Some immunomodulatory properties that have been reported include enhancement of host resistance against infection, anti-inflammatory effect, inhibition of tumor growth, as well as modulation of some basic immune function at the cellular level. American ginseng, Panax quinquefolium, is another specie of ginseng which has gained popularity as a health supplement having many beneficial health effects. Several groups of scientists have attempted to isolate and elucidate the structure of the polysaccharides present in ginseng. Some of the polysaccharides have been demonstrated to be active in modulating the immune system.\nA series of studies on the isolation, characterization, and biological evaluation of ginseng polysaccharides was carried out by Tomoda's group in Kyoritsu College of Pharmacy, Japan. In one set of studies, ginseng polysaccharides were fractionated based on their acidity. Two acidic polysaccharides having immunological activities have been isolated from root of Korean ginseng (Panax ginseng)[1,2]. The sliced roots were extracted with hot water. The extract was treated with cetyltrimethylammonium bromide (CTAB) in the presence of sodium sulfate. The precipitate was separated, dialyzed, and applied to a Sephadex G-25 column, DEAE-Sephacel (Pharmacia) column to give two pure polysaccharides, which were designated as ginsenan PA and ginsenan PB. Gel chromatography on Toyopearl HW-55F gave the values of 1.6×105 and 5.5×104 for the molecular weight of ginsenan PA and ginsenan PB, respectively. Quantitative analyses showed that ginsenan PA contained 21.3% arabinose, 53.4% galactose, 2.0% rhamnose, 16.0% galacturonic acid and 2.7% glucuronic acid. The molar ratio of these component sugars was 11:22:1:6:1. Ginsenan PB contained 11.0% arabinose, 32.2% galactose, 8.1% rhamnose, 39.9% galacturonic acid, and 5.0% glucuronic acid. The molar ratio was 3:7:2:8:1. Both polysaccharides showed marked reticuloendothelial system-potentiating activity in a carbon clearance test, and pronounced anti-complementary activity and alkaline phosphatase-inducing activity in a dose dependent manner.\nIn another study[3], an additional two polysaccharides were isolated from the supernatant of the above extract treated with CTAB, i.e., the supernatant was poured into ethanol. The precipitate was separated and applied to columns of DEAE-Sephadex A-25 and Sephadex G-25 to give another two pure polysaccharides, designated S-IA and S-IIA. Gel chromatography on Toyopearl HW-55F gave the values of 5.6×104 and 1.0×105 for the molecular weight of S-IA and S-IIA respectively. Ginsenan S-IA contains 42.3% arabinose, 50.8% galactose and 6.9% galacturonic acid with the molar ratio of 8:8:1. Ginsenan S-IIA is composed of 42.0% L-arabinose, 32.6% galactose, 6.2% glucose, and 19.2% galacturonic acid. The molar ratio is 15:10:2:5.\nSeveral polysaccharide fractions from leaves and roots of Panax ginseng have been separated by Yamada's group in the Oriental Medicine Research Center of the Kitasato Institute, Japan. The chemical properties and biological activities were investigated and compared.[4]\nGinseng roots and leaves from China, after treated with ethanol to remove their ginsenosides, were extracted with water, and the residue were extracted with 0.5 M NaOH to give water-soluble and alkaline soluble polysaccharide fractions (designated as GR-2 and GL-2 for the water-soluble fractions and GRA-2 and GLA-2 for the alkaline-soluble fractions, respectively). Based on the acidity of the component polysaccharides, all fractions were further fractionated into strongly acidic (designated as GR-3, GL-3, GRA-3, and GLA-3), weakly acidic (designated as GR-4, GL-4, GRA-4, and GLA-4), and neutral polysaccharide fractions (designated as GR-5, GL-5, GRA-5, and GLA-5), by treatment with cetyltrimethylammonium. The roots contained larger amount of polysaccharide than the leaves. The strongly acidic polysaccharide fractions from the roots had a high content of uronic acid, even higher than 50%. Similar component sugars were detected from all fractions. They were rhamnose, arabinose, galactose, glucose, galacturonic acid, and glucuronic acid. Galacturonic acid was the main uronic acid component.\nThe fraction with highest anti-complementary activity, GL-3, was further fractionated with columns of DEAE-Sephadex, Sepharose CL-6B, DEAE Toyopearl, and ethanol precipitation to give fractions designated as GL-PI through GL-PIV.[5]All fractions contained 32-44% uronic acid. Fraction PI had the highest molecular weight of 50,000. PI and PII consisted mainly of Rha, Gal, and GalA, and PIII contained Fuc in addition, whereas PIV consisted of Gal, Glc, and GalA. A detailed structural determination was performed.\nAn anti-ulcer pectic polysaccharide (GL-BIII) was isolated from weakly acidic polysaccharide fraction GL-4 by chromatography on DEAE Sepharose CL-6B and Sepharose CL-6B. It was mainly composed of Rha, Ara, Man, Gal, Glc, GalA, and GlcA in the molar ratio of 3:4:2:10:1:7:4. Detailed structural determination was performed.[6]\nAnother macrophage Fc receptor expression-enhancing polysaccharide (GL-4IIb2) was separated from GL-4 by anion-exchange chromatography on DEAE-Sepharose CL-6B. Chemical analysis showed that the sample contained 65% total carbohydrate and 33.7% uronic acid. The composition analysis and structural determination were performed.[7]\nAnother Panax ginseng extract with anticomplementary activity, G-115, was studied by the same group.[8]G-115 was fractionated in order to characterize the active substances for anticomplementary and mitogenic activities. The most potent anticomplementary activity was observed in the crude polysaccharide fraction, G-115G, whereas the water-soluble dialyzable fraction, G-115E, showed the most potent mitogenic activity. G-115G was further purified by precipitation with cetyltrimethylammonium bromide, anion-exchange chromatography on DEAE-Sepharose and gel filtration on Sepharose CL-4B, and a major potent anticomplementary polysaccharide, G-115I1-IIa-2-3 was obtained. This polysaccharide was homogeneous. Its molecular weight was estimated to be 3.68×105. It consisted mainly of arabinose, galactose and glucose in addition to small amounts of galacturonic acid, glucuronic acid and rhamnose.\nGinseng polysaccharides were isolated from Korean, Chinese, and Japanese ginseng by Hikino's group at the Pharmaceutical Institute, Tohoku University, Japan. The hypoglycemic activities of the ginseng polysaccharides have been tested. The composition and some structure feature have been elucidated.[9-14]\nThree polysaccharides, quinquefolans A through C, were isolated from American ginseng.[14]Their molecular weights were estimated to be higher than 2.0×106 by gel chromatography over Sephacryl S-500. The neutral sugar components were mannose and glucose (molar ratio, 1.0:2.3) for quinquefolan A, mannose and glucose (1.0:5.5) for quinquefolan B, and xylose for quinquefolan C. The acidic sugar components in quinquefolans A through C were found to be 10.8, 11.7, and 7.1% respectively. The content of peptide moieties in these glycans was 2.7, 2.9, and 2.3% for quinquefolans A through C respectively. All of these polysaccharides showed hypoglycemic effects in normal and alloxan-reduced hypoglycemic mice.\nAn acidic polysaccharide with the molecular weight of 150,000, called ginsan, was isolated from Panax ginseng by a research group at the Laboratory of Immunology, Korean Cancer Center Hospital, Seoul, Korea.[15]This polysaccharide was composed of 3.7% protein and 47.1% hexose (glucose and galactose) and 43.1% uronic acid (galacturonic acid). Ginsan induced the proliferation of T cells and B cells and generated lymphokine activated killer cells from both natural killer and T cells through endogenously produced multiple cytokines.[16]\nMiao et al. from Northeast Normal University of China has isolated polysaccharides from American ginseng. The purification and structural analysis were performed.[17]\nThe biological activities of polysaccharides from American ginseng have been investigated by a research group in Norman Bethune University of Medical Science.[18, 19]They found that polysaccharides from American ginseng (PPQ) enhanced lymphocyte transformation. The effect of polysaccharide from Panax quinquefolium (PPQ-1) on cytokine production from murine spleen lymphocyte in vitro was studied. The data suggest that PPQ-1 regulates immune function."}
{"text": "This invention relates generally to electromagnetic relays and more specifically to flat electromagnetic relays which have a very small width when vertically oriented and or a very small height when horizontally oriented.\nFlat electromaqnetic relays are already known from U.S. Pat. No. 4,010,433 granted to Hiromi Nishimura et al Mar. 1, 1977; U.S. Pat. No. 4,031,493 granted to Michael Van Der Wielen Jun. 21, 1977; U.S. Pat. No. 4,272,745 to Takashi Tanaka Jun. 9, 1981; U.S. Pat. No. 4,290,037 granted to Takashi Inagawa et al Sep. 15, 1981; U.S. Pat. No. 4,517,537 granted to Josef Weiser et al May 14, 1985 and from U.S. Pat. No. 4,684,909 granted to Michael Dittmann Aug. 4, 1987.\nA flat electromagnetic relay is also already known from U.S. Pat. No. 5,038,123 granted to Christopher Alan Brandon Aug. 6, 1991. This patent which is incorporated herein by reference discloses a flat electromagnetic relay comprising a lead assembly having six stamped insert molded circuit leads, an armature frame pivotally supporting a balanced beam armature which carries two contact bars, a coil assembly and an electromagnetic frame having two diagonally arranged pole wings which are on opposite sides and ends of the armature. The armature is biased into a first operative position where the contact bar at one shunts two circuit leads. When the coil is energized, the armature is pivoted into a second operative position where the contact bar at the other end shunts two other circuit leads.\nWhile this flat electromagnetic relay has many advantages it has a disadvantage in that the armature assembly is complicated, expensive to manufacture and does not lend itself to miniaturization very well."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of Sic Bo dice games, both physical and electronic driven dice games, to variations of wagers for Sic Bo variants and provision of physical and electronic variations in random selection of dice used in gaming outcomes.\n2. Background of the Art\nSic Bo, meaning “dice pair” is an ancient Chinese gambling game. It is very popular in Macau, in my estimation second to baccarat only. In the United States is often found in the Asian gaming rooms, especially in Atlantic City. The large Las Vegas properties will usually have one table in the general casino.\nThe game uses three dice and a table with a variety of betting options on the roll of those dice. The odds and table layout may also vary from place to place. However, the payoffs are the same across Atlantic City and Macau, but different from each other. In Vegas and on the Internet casinos, anything is possible.\nFollowing is a list of the bets available. The payoffs vary on some bets, from casino to casino, so for those bets a range of viable payoffs is indicated.                Small: Wins on total of 4-10, except for a three of a kind. Probability of winning is 48.61%. Pays 1 to 1. House edge is 2.78%.        Big: Wins on total of 11-17, except for a three of a kind. Probability of winning is 48.61%. Pays 1 to 1. House edge is 2.78%.        Total of 4 and 17: Pays 50 to 1 in Macau, 60 to 1 in Atlantic City, and 62 to 1 in Australia.        \nHouse Edge for 4 and 17 BetsPaysProbabilityReturn501.39%29.17%511.39%27.78%521.39%26.39%531.39%25.00%541.39%23.61%551.39%22.22%561.39%20.83%571.39%19.44%581.39%18.06%591.39%16.67%601.39%15.28%611.39%13.89%621.39%12.5%631.39%11.11%641.39%9.72%651.39%8.33%                Total of 5 and 16: Pays 18 to 1 in Macau, 30 to 1 in Atlantic City, and 31 to 1 in Australia.        \nHouse Edge for 5 and 16 BetsPaysProbabilityReturn182.78%47.22%192.78%44.44%202.78%41.67%212.78%38.89%222.78%36.11%232.78%33.33%242.78%30.56%252.78%27.78%262.78%25.00%272.78%22.22%282.78%19.44%292.78%16.67%302.78%13.89%312.78%11.11%322.78%8.33%                Total of 6 and 15: Pays 14 to 1 in Macau, 17 to 1 in Atlantic City, and 18 to 1 in Australia.        \nHouse Edge for 6 and 15 BetsPaysProbabilityReturn144.63%30.56%154.63%25.93%164.63%21.3%174.63%16.67%184.63%12.04%194.63%7.41%                Total of 7 and 14: Pays 12 to 1 everywhere.        \nHouse Edge for 7 and 14 BetsPaysProbabilityReturn106.94%23.61%116.94%16.67%126.94%9.72%                Total of 8 and 13: Pays 8 to 1 everywhere.        \nHouse Edge for 8 and 13 BetsPaysProbabilityReturn79.72%22.22%89.72%12.50%                Total of 9 and 12: Pays 6 to 1 in Macau and Atlantic City. Pays 7 to 1 in Australia.        \nHouse Edge for 9 and 12 BetsPaysProbabilityReturn611.57%18.98%711.57%7.41%                Total of 10 and 11: Pays 6 to 1 everywhere. House edge of 12.50%.        Double: Bet on any specific pair (for example two l's). Player wins if at least 2 of the number chosen appears. Pays 8 to 1 in Macau, 10 to 1 in Atlantic City, and 11 to 1 in Australia.        \nHouse Edge for Double BetsPaysProbabilityReturn87.41%33.33%97.41%25.93%107.41%18.52%117.41%11.11%                Triple: Player may bet on any specific trips (for example three l's). Player wins if all 3 dice match the number chosen. Pays 150 to 1 in Macau and 180 to 1 in Atlantic City.        \nHouse Edge for Triple BetsPaysProbabilityReturn1500.46%30.09%1550.46%27.78%1600.46%25.46%1650.46%23.15%1700.46%20.83%1750.46%18.52%1800.46%16.20%1850.46%13.89%1900.46%11.57%                Any Triple: Wins on any three of a kind. Pays 24 to 1 in Macau and 30 to 1 in Atlantic City.        \nHouse Edge for Any Triple BetPaysProbabilityReturn242.78%30.56%252.78%27.78%262.78%25.00%272.78%22.22%282.78%19.44%292.78%16.67%302.78%13.89%312.78%11.11%322.78%8.33%                Any One Number: Bet on any specific number from 1 to 6. If chosen number appears 1 time bet pays 1 to 1, if it appears 2 times bet pays 2 to 1, and if it appears 3 times it pays 3 to 1 (in Atlantic City and Macau). Probability of 1 match is 34.72%, 2 matches is 6.94%, 3 matches is 0.46%. House edge of 7.87% (in Atlantic City and Macau). In Australia three matches pays 12 to 1, for a house edge of 3.70%.        Odd: Wins on any odd total, except loses with any three of a kind. The probability of winning is 48.61% and house edge is 2.78%.        Even: Wins on any even total, except loses with any three of a kind. This bet is predominantly in Macau, not in the United States. Probability of winning is 48.61% and house edge is 2.78%.        xyz: The player chooses a three-number combination, using three unique numbers, for example 1-3-5. The bet pays 30 to 1 if all three are rolled. Probability of winning is 2.78% and house edge is 13.89%.        xxy: The player chooses a three-number combination, using two unique numbers, and specifying which will be rolled twice. For example 2-2-4. The bet pays 50 to 1 if the roll matches exactly. So, in this example, 2-4-4 and 2-4-5 would lose. Probability of winning is 1.39% and house edge is 29.17%.        wxyz: The player chooses a four-number combination, using four unique numbers. For example 1-3-4-6. The bet pays 7 to 1 if the roll matches three of the chosen numbers. In my example, 1-3-4, 1-3-6, 1-4-6, and 3-4-6 would be the only winners. Probability of winning is 11.11% and house edge is also 11.11%.        xy: The player chooses a two-number combination, using two unique numbers. For example 2-5. The bet pays 5 to 1 if the roll contains both numbers. For example a bet on 2-5 would win on a rolls of 1-2-5, 2-2-5, 2-3-5, 2-4-5, 2-5-5, and 2-5-6. Probability of winning is 13.89% and house edge is 16.67%.        \nFollowing are tables summarizing each bet under the Atlantic City and Macau rules.\nSic Bo - Atlantic City RulesWinningBetPaysCombinationsProbabilityReturnSmall, Big11050.486111−0.0277784, 176030.013889−0.1527785, 163060.027778−0.1388896, 1517100.046296−0.1666677, 1412150.069444−0.0972228, 138210.097222−0.1259, 126250.115741−0.18981510, 11 6270.125000−0.125000Triple18010.004630−0.162037Any triple3060.027778−0.138889Double10160.074074−0.185185Xy5300.138889−0.166667Any number1, 2, 375, 15, 10.421296−0.078704\nFollowing is a table summarizing each bet under the Macau rules.\nSic Bo - Macau RulesWinningBetPaysCombinationsProbabilityReturnSmall, Big11050.486111−0.027778Odd, Even11050.486111−0.0277784, 175030.013889−0.2916675, 161860.027778−0.4722226, 1514100.046296−0.3055567, 1412150.069444−0.0972228, 138210.097222−0.1250009, 126250.115741−0.18981510, 11 6270.125000−0.125000Triple15010.004630−0.300926Any triple2460.027778−0.305556Double8160.074074−0.333333xyz3060.027778−0.138889xxy5030.013889−0.291667wxyz7240.111111−0.111111xy5300.138889−0.166667Any number1, 2, 375, 15, 10.421296−0.078704Sic Bo - Australia RulesWinningBetPaysCombinationsProbabilityReturnSmall, Big11050.486111−0.0277784, 176230.013889−0.1250005, 163160.027778−0.1111116, 1518100.046296−0.1203707, 1412150.069444−0.0972228, 138210.097222−0.1250009, 127250.115741−0.07407410, 11 6270.125000−0.125000Triple18010.004630−0.162037Any triple3160.027778−0.111111Double11160.074074−0.111111Xy6300.138889−0.027778Any number1, 2, 1275, 15, 10.421296−0.037037\nFollowing is a table summarizing each bet under the Australia rules.\nSic Bo House Edge ComparisonBetMacauAtlantic CityAustraliaSmall, Big2.78%2.78%2.78%4, 1729.17%15.28%12.50%5, 1647.22%13.89%11.11%6, 1530.56%16.67%12.04%7, 149.72%9.72%9.72%8, 1312.50%12.50%12.50%9, 1218.98%18.98%7.41%10, 11 12.5%12.50%12.50%Triple30.09%16.20%16.20%Any triple30.56%13.89%11.11%Double33.33%18.52%11.11%xy16.67%16.67%2.78%Any number7.87%7.87%3.7%Yee Hah Hi\nAbout 20% of the many Sic Bo tables in Macau use dice with pictures, instead of numbers. The name of this game is “Yee Hah Hi.” The translation is “Fish Shrimp Crab.” Each picture has an associated number, as shown below. In addition to the usual bets there are some based on the colors of the symbols, as explained below.\nSic Bo - Chinese Dice RulesSymbolColorNumberFishRed1ScorpionGreen2GoardBlue3CoinBlue4CrabGreen5RoosterRed6                All same color: Wins if all three dice are the same color. Pays 7 to 1. The probability of winning and house edge are both, coincidentally, 11.11%.        Exactly two of specific color: Wins if exactly two dice are a color chose by the player, for example red. Pays 3 to 1. The probability of winning is 22.22%, and house edge is 11.11%.        Exactly one of specific color: Wins if exactly one die is a color chose by the player, for example red. Pays 1 to 1. The probability of winning is 44.44%, and house edge is 11.11%.Underlying Rules of Sic BoSetup        \nSic Bo is a dice game which consists of many places on the table to place bets. Three dice are used to determine the outcome of the betting.\nBetting\nThe bottom row of the Sic Bo table has six betting areas. These are called single number bets. If one of the three dice comes up as one the numbers bet on (from 1-6), you get paid out at 1 to 1. If two numbers come up you get paid out at 2 to 1 and if three numbers come up you get paid out at 3 to 1. So if you bet on the number 3 and you wagered $5.00 and 2 threes came up you would win $10.00.\nThe two number combination bets consist of placing a wager on any 2 dice combinations. If you wager on a 6,4 combination and the result of the dice roll was 4, 7, 6, you would win at 6 to 1 odds. You can only win on one instance of the two-number combination.\nA three number total wager consists of combining the total shown on the resulting dice roll. Different totals have different odds. A result of 3 or 18 is always a loss. The payouts are as follows:\n4 or 17, 50 to 1;\n5 or 16, 25 to 1;\n6 or 15, 15 to 1;\n7 or 14, 10 to 1;\n8 or 13, 6 to 1;\n9 or 12, 5 to 1;\n10 or 11, 5 to 1.\nYou can wager on whether the combined total of the dice roll will be either between 4-10 or 11-17. This wager pays out at 1-1 odds. All small or big wagers lose if the result of the dice is a triplet.\nTriplets or trips (Trips) are when the dice roll is all one number. So if you wagered on three 6's and that was the dice result you would win at 150-1 odds. You can also wager on all six triplet bets at once, this is called “Any Triplet” and pays out at 25-1 odds.\nPairs are when the resulting dice roll contains a pair of the same number. Whether you have two fives or three fives you still only win once and you win at odds of 11-1.\nSic Bo is a single roll dice game in which players may wager on the sum of the dice, typically three. For example, the player may wager on the sum of twelve without regard to the combination of the dice required to add up to twelve, i.e. the player would win if the dice show 6-5-1, or 5-5-2. Alternatively or additionally, the player may wager on the combination of the dice. In the before-mentioned example, the player might wager on the roll of the dice results to be 6-6-6. In Sic Bo, three dice are rolled and all wagers are immediately resolved; that is, there are, at present, no multiple roll wagers in Sic Bo. Sic Bo is an easy dice game to understand and play, however its simplicity fails to keep the player interested and large payouts can be limited.\nSic Bo, known to some people as Tai Sai, means dice pair in Chinese. It is an ancient Chinese game of chance. It may have originally used dice shaken between a plate and an overturned bowl, but today it commonly has three dice in a cage for tumbling. Extremely popular in Asian cultures, this exciting and engaging game is winning new converts in many casinos around the world. Sic Bo is easy to play. The object is to pick the numbers or combinations that will appear on the dice when they come to rest after tumbling. There are 50 different ways to place bets, so players have plenty of choices with a range of payouts—some as high as 180-to-1. In a sense, this game is somewhat similar to the Western game of craps, but it involves many more betting opportunities and combinations due to the use of one additional dice than the game of craps. However, despite this similarity, this game has not yet found its way into traditional Western-style casinos.\nA dice game proposed in U.S. Pat. No. 5,413,351 as described above is performed according to special rules, however, there is a problem that people are not familiar with this game and its special rules and therefore this game lacks amusement. The Sic Bo generally widely known is familiar to people, so that it can amuse players, however, the BET patterns are limited, so that there is still room for upgrading in terms of improvement in amusement. In detail, in the BET region whose appearance frequency is low and award ratio is highest (approximately 1 to 180), it is considered that a player performs a BET operation in an excited state, however, in this BET region, it is only predicted that rolled numbers on three dice rolled are the same (the combination of rolled numbers on the dice is (1, 1, 1), (2, 2, 2) . . . (6, 6, 6)), and this area lacks excitement when making a BET.\nVariations of Sic Bo are disclosed as Published US Patent Application Document No. 20060049578 (Chen) describes a game titled Let Me Roll Sic Bo as a dice game blending both Sic Bo with the traditional craps game. There are three regular dice (cubes). Each die has six sides and different number of spots on it (from one to six). Each player takes a turn to shoot the dice clockwise. Players can wager on nine main designed betting areas to win different odds.\nAnother variation of Sic Bo are disclosed as Published US Patent Application Document No. 20060202416 (Madden) describes a method of playing a betting dice game with a progressive payout is disclosed. The game includes three dice and a game table with up to four wager zones for players to place a bet on. The wager zones correspond to a sum of the dice. Payout to a winning player results from the wager placed on the wager zones and the corresponding first, second and third sequential rolls of the dice.\nPublished US Patent Application Document No. 20060249907 (Wong) describes an East-West casino game that offers two or more gambling games in which at least one game involves a Feng Shui element of chance and the other involves either a game based upon a new Chinese Poker Deck of cards or a video slot machine having a carousel in which two or more wheels rotate in opposite directions and a random event associated with the two or more wheels is used in connection with a payout or bonus round of the video slot machine. The game involving a Feng Shui element of chance can be a Yin-Yang roulette game or a video slot machine (which may at least partially based upon the Chinese Poker Deck) or a Feng Shui keno machine while a game based upon the Chinese Poker Deck may be a poker game and the casino itself may be a website. The video slot machine can use a carousel in which two or more wheels (which need not be round) rotate in opposite directions and additional wheel can be added in which symbols are shuffled. Examples of video slot machines include a Yin-Yang roulette game in which two balls travel in opposite directions (similar to a roulette wheel) with landings for the ball designated by symbols related to the Chinese games of Tin Gau, Mahjongg, Fan Tan or Sic Bo and a Feng Shui keno machine in which at least one Chinese character is broken down into a plurality of radicals.\nPublished US Patent Application Document No. 20100069142 (Kido) describes a dice gaming machine accepts a bet on a bet area and a bet on an intermediate area between/among a plurality of bet areas. After ending the bet acceptance, the gaming machine starts rolling dice. After the rolling of the dice has been stopped, the dice gaming machine identifies a winning bet area based on the dice rolling result. A payout amount for each terminal is calculated based on the identified winning bet area, the bet area or the intermediate area on which the bet is accepted, and a payout magnification ratio set for the bet area or the intermediate area.\nAutomatic Sic Bo Machines are known in the art wherein dice are automatically tossed and read, and wagers are electronically entered and resolved based on the automatic reading and tossing of dice. Combination systems where manually tossed dice, manually entered outcomes, and electronic wagers are used that are resolved based upon the manually entered (or electronically read) manually thrown dice outcomes. TCS John Huxley showcased a machine at 2007 G2E Asia in Macau. The ergonomically designed Super Sic Bo table featured a high definition 55 inch LCD playing screen with an apparent Sic Bo playing surface. High impact computer generated graphic animations highlighted ‘no more bets’ and all winning bets as well as indicating the streak bets of either odd or even and big or small numbers as they progress on the Super Streak side bet.\nA voice announcement feature delivered game results and number sequences. This was available in multiple languages and gave players concise and distinctive game information and attracted new players to join in.\nThere was also offered a Super Streak side bet option offering players the opportunity to participate in a bonus bet determined by the results over a number of games, which are clearly displayed and easy to understand. This provides incremental payouts for a streak of either odd or even and big or small results. A triple result where the three dice all display the same number does not cancel the Super Streak.\nThe system was compatible with e-FX LCD displays and with Animator LED displays, showing the last 17 individual dice results, big or small betting options, odd or even betting options and Super Streak side bet streak progress are easy to see.\nSystem Highlights Included                Jumbo size color dices (888 mm)        Integrated electronic dice shaker        Electronic dice outcome recognition        Digital dice results communication        Touch-screen bet placing        Fully automated game        SAS 6.02 enabled        Adaptive design and graphic presentation        Minimum and maximum bet limits        Wide range of betting options and payout odds        Automatic game results calculation and payout        Increased productivity—no game stops        Jackpots and side bets available        Betting time indication        Easier monitoring and accounting        Game information display        Electronically activated        Precise dice rolling        Elimination of “no spin” and incorrect dice throw        Special system for flat lying dices        Dice shaker failure detection        Glass dome-shaped shaker cover        Digital Dice recognition        Touch-screen betting        No human interference        Faster game playable 24/7 without personnel in attendance        Game and dice statisticsAn example image of an electronic system is shown in FIG. 1. All references cited herein are incorporated herein by reference.        "}
{"text": "An electrical power system operates under a steady-state condition when there exists a balance between generated and consumed active power for the system. Power system disturbances may cause oscillations in machine rotor angles that can result in conditions like a power swing, when internal voltages of system generators slip relative to each other. Power system faults, line switching, generator disconnection, and the loss or sudden application of large amounts of load are examples of system disturbances that may cause the power swing event to occur in the system. Depending on the severity of the disturbance and power system control actions, the system may return to a stable state or experience a large separation of load angle and eventually lose synchronism. Large power swings, stable or unstable, may cause unwanted relay operations at different locations in the system, which can aggravate the system disturbance and can result in major power outages or blackouts.\nFurther, asynchronous operation of interconnected generators in the power system as an effect of unstable power swing may initiate uncontrolled tripping of circuit breakers resulting in equipment damage and posing a safety concern for utility operators. Therefore, the asynchronous system areas may need to be separated from each other quickly and dynamically in order to avoid extensive equipment damage and shutdown of major portions of the system. In order to contain these risks, it is required as per international standards to have an optimal generator protection device such as a generator relay in place to isolate generators from rest of the system within a half-slip cycle. The need to meet the international standards challenges protection engineers to ensure selective and reliable relay operation.\nIn a conventional relaying approach, a variation in system impedance determined at generator terminals is analyzed for detecting power swing. Various impedance-based protection approaches including power swing block (PSB) and out-of step trip (OST) are being currently used. However, these protection approaches may need an extensive power system stability study to arrive at an optimal setting for selective and reliable relay operation. Protection engineers typically use preliminary settings that are not adapted to accommodate variation in system configurations or operational dynamics, for example, changes in transmission and distribution layout during implementation phase or dynamically during operational phase. Extensive study and non-dynamic preliminary settings may result in the protection device being unable to selectively, reliably and dependably detect unstable power swings and isolate generators during such events.\nOther known relaying approaches estimate swing center voltage (SCV) for detecting power swings. Such approaches use approximate estimation that does not take into consideration real time power system dynamics. In some relaying approaches, a high-speed communication network such as fiber optic or global positioning system (GPS) communication is used to obtain data at a source end from one or more generators at receiving end(s), which is at a remote location from the source end, for SCV estimation. However, such approaches have economic challenges due to cost associated with implementing and maintaining high-speed communication network. Moreover, such approaches may require data to be determined at the receiving end and then time synchronize the receiving end data with data determined at the source end."}
{"text": "1. Field of the Invention\nThis invention relates to a dielectric ceramic composition for a temperature-compensating ceramic capacitor which has a high dielectric constant and in which the temperature coefficient of its capacity can be controlled within the range of .+-.30 ppm (10.sup.-6)/.degree.C. at temperatures within the range of -55.degree. to +125.degree. C. The invention also relates to a multilayer ceramic capacitor obtained by firing at relatively low temperatures a laminate of a layer of the above dielectric ceramic composition and an Ag-Pd alloy layer containing a major amount of silver.\n2. Description of the Prior Art\nGenerally, a temperature-compensating ceramic capacitor having a high dielectric constant, particularly a commercial multilayer ceramic capacitor, is obtained by laminating a plurality of thin dielectric layers having an inner electrode formed on their surface so that the internal electrodes are alternately connected in series to electrodes for external connection, and firing the assembly. Laminated capacitors of this type require relatively high firing temperatures (at least 1240.degree. C.) in order to make them sufficiently compact and thus give a high dielectric constant. Hence, the electrodes must be made of a metal having a higher melting point than the firing temperature of the dielectric material, such as an expensive noble metal (e.g., platinum or palladium). The cost of the metallic material thus increases the total cost of capacitors of this type.\nAttempts have therefore been made to obtain multilayer ceramic capacitors of low prices by lowering the firing temperature for the dielectric material and the internal electrodes, and using inexpensive Pd-Ag alloys particularly having a major amount of Ag instead of the expensive noble metals. It is known however that generally, the crystallinity of a dielectric decreases with decreasing firing temperature, and therefore, its dielectric constant decreases. When the firing is carried out at lower temperatures than the aforesaid firing temperature, sufficient electrical properties and temperature characteristics for a high dielectric ceramic capacitor cannot be obtained.\nJapanese Laid-Open Patent Publication No. 170405/1982 states that by adding moderate amounts of ZnO and SiO.sub.2 to a composition comprising Nd.sub.2 TiO.sub.7, BaTiO.sub.3, TiO.sub.2, Bi.sub.2 O.sub.3 and Pb.sub.3 O.sub.4, the resulting composition can be fired at 1050.degree. to 1100.degree. C. to obtain a sintered body, and in the production of a laminated capacitor in which the baking of the internal electrodes and the firing of the ceramic composition must be simultaneously carried out, a low-melting inexpensive Ag-Pd alloy can be used as the internal electrodes. This prior technique, however, has the defect that the temperature coefficient of the electrostatic capacity is within a relatively broad range of +30 to -300 ppm/.degree.C., the insulation resistance is reduced, or the breakdown voltage is low and varies greatly."}
{"text": "Lithium ion secondary batteries are widely used for portable electronic instruments such as mobile phones or notebook-size personal computers. As a cathode active material for a lithium ion secondary battery, a composite oxide of lithium with a transition metal, etc. (hereinafter sometimes referred to as lithium-containing composite oxide) such as LiCoO2, LiNiO2, LiNi0.8Co0.2 O2 or LiMn2O4, is employed. In recent years, it is desired to further reduce the size and weight as a lithium ion secondary battery for portable electronic instruments or vehicles, and a further improvement in the discharge capacity per unit mass and in the characteristics such that the discharge capacity will not be decreased after charge and discharge cycles are repeatedly carried out (hereinafter sometimes referred to as cycle characteristic), is desired.\nTo improve the cycle characteristics, it has been known to be effective to form a covering layer on a lithium-containing composite oxide.\nPatent Document 1 discloses a process of adding an ammonium fluoride aqueous solution to a dispersion having a lithium-containing composite oxide dispersed in an aluminum nitrate aqueous solution, followed by filtration, washing and heating to form a covering layer of aluminum fluoride on the surface of the lithium-containing composite oxide. However, by this process, a step of carrying out filtration and washing is essential, whereby the process is complicated and in addition, a waste disposal treatment is necessary, and thus the process is inferior in the productivity. Further, when the wet cake after filtration is dried, the cathode active material tends to be agglomerated to form coarse particles.\nPatent Document 2 discloses a method of fluorinating a cathode material using a fluorinating agent such as fluorine (F2), nitrogen trifluoride or chlorine trifluoride. However, this method requires special equipment since a highly toxic gas is used.\nPatent Document 3 discloses a process of stirring and mixing an aqueous solution containing zirconium and a lithium-containing composite oxide represented by the formula LipNxMyOzFa (0.9≦p≦1.1, 0.965≦x<1.00, 0<y≦0.035, 1.9≦z≦2.1, x+y=1 and 0≦a≦0.02) wherein the molar amount of Li element is from 0.9 to 1.1 molar times the total molar amount of the transition metal element, and firing the mixture at high temperature of at least 450° C. in an oxygen atmosphere to obtain a cathode active material having the surface layer of the lithium-containing composite oxide covered with zirconium oxide. By this method, it is difficult to cover the lithium-containing composite oxide with a compound other than an oxide."}
{"text": "Subsurface resources such as oil, gas and water are typically recovered by drilling a well bore from the surface to a subterranean reservoir or zone that contains the resources. The well bore allows oil, gas and water to flow to the surface under its own pressure. For low pressure or depleted zones, rod pumps are often used to retrieve the fluids to the surface.\nTo facilitate drilling and production operations, cavities are sometimes formed in the production zone. Short extensions, or “rat holes,” are often formed at the bottom of the cavity to collect cuttings and other drilling debris. As the subsurface liquids collect in the well bore, the heavier debris falls to the bottom of the rat hole and is thereby both centralized and collected out of the cavity. To avoid being clogged with debris, a pump inlet may be positioned within the cavity above the rat hole. The pump inlet may be positioned fairly low in the cavity (for example, below the fluid waterline) to avoid vapor lock. Traditional methods of positioning a pump inlet are sometimes inaccurate and inefficient, leading to clogging or vapor lock and increased maintenance and operations costs for the well."}
{"text": "1. Field of the Invention\nThe present invention relates to a printing control apparatus and a printing control method used to execute printing using recording material stored in a plurality of storage portions.\n2. Description of the Related Art\nIt is known that, when a conventional printing apparatus having a plurality of storage portions capable of storing recording material such as recording sheets executes printing, such printing apparatus uses the recording material stored in the storage portions while switching the storage portions. With such a printing apparatus, the user checks the type of recording material used by a print job to be executed and sets appropriate recording material in advance. According to a technique discussed in Japanese Patent Application Laid-Open No. 2005-193486, even after printing is started, the use order of a plurality of storage portions storing sheets identical to the sheet that is being used and storage portions into which the currently used sheet can be supplied are indicated. Namely, the technique enables supply of a sheet even during a print process.\nHowever, the technique discussed in Japanese Patent Application Laid-Open No. 2005-193486 assumes that only the cut sheets are used as recording material. Use of a continuous sheet, such as a roll sheet, as recording material is not assumed. When a continuous sheet is used as recording material, unlike cut sheets, the continuous sheet cannot be supplied by addition. The recording material cannot be used appropriately by simply presenting unused storage portions."}
{"text": "The present disclosure relates to electronic devices capable of apparatus setting corresponding to a region where it is installed, storage media, and apparatus setting methods.\nAmong image forming apparatuses such as digital multifunction peripherals and printers, a certain image forming apparatus identifies its self-location using a satellite positioning system such as global positioning system (GPS), and performs apparatus setting on a regional basis according to a local resource information item corresponding to the self-location.\nFor example, a certain printer includes therein a GPS to identify its self-location at a power-up using the GPS. The printer then downloads a font (local resource information item) necessary in the region that includes the self-location from a font server.\nThis can enable appropriate setting of the printer corresponding to the region where the pointer is located. However, the necessity of an external server such as the font server may incur an increase in cost.\nIn view of the foregoing, another printer stores therein fonts in several languages in a switchable manner to eliminate an external server.\nHowever, a license fee may be required for use of some font. Accordingly, for storage of every font in a switchable manner in the printer, an expense for the font requiring the license fee may be incurred regardless of the use or non-use of the font. This may limit cost reduction."}
{"text": "Protein receptors are known normally to bind to their target ligands via epitopes, which constitute a small proportion of the total protein molecule. For maximum binding or interaction, the structure of the epitope needs to be maintained in a rigid conformation in order to form a binding site containing all the necessary components of the epitope in close proximity. Attempts to produce an analogous peptide constructed solely of the amino acids comprising the binding site often fail because these peptides do not possess the same biological activity as the protein receptor. This is attributed to the peptide having a different conformation in free solution from that of the entire protein receptor. In addition, where the binding site of a protein is constructed of oligo-peptides from different, non-contiguous parts of a protein chain, mixing isolated oligopeptides in free solution does not result in reconstitution of the active binding site.\nBeing constrained to use such large proteins to present binding-site epitopes gives rise to several problems in development of new receptor-specific therapeutic strategies. One problem is that such large proteins can readily evoke an immune response. A second problem is that long peptide chains are susceptible to attack by endopeptidases, such as those in the lumen of the gut. Finally, these large proteins can be costly to manufacture, purify and maintain in stable form."}
{"text": "In-vehicle entertainment has traditionally been provided by the radio. The driver and, if present, passengers in the vehicle usually select a radio station, and then passively listen to the audio, be it music or talk. Other audio entertainment is also available, including cassette tapes, CDs, audiobooks, portable electronic music players. In each case, the vehicle occupants passively consume the content presented to them aurally.\nMore recently, in-vehicle entertainment systems provide the opportunity for vehicle passengers to watch video, such as movies or TV programs, on embedded in-vehicle displays. Such an option is not typically available to the driver of the vehicle due to the likelihood of distraction. And again, such entertainment options consist of the vehicle occupants passively consuming the content of the movie or TV program.\nMore advanced entertainment systems may include gaming consoles or computers. These systems give the vehicle passengers the opportunity to engage in, for example, video games while traveling to a destination. Again, however, such systems are not typically available to the driver of the vehicle because of the concern that such a game system could be distracting to the driver."}
{"text": "1. Field of the Invention\nThe present invention relates to a driving force transmission device made of synthetic resin.\n2. Description of the Related Art\nBecause gears made of synthetic resin can be readily formed by injection molding or the like at a low cost, manufacturing costs of drives employing synthetic resin gears can be reduced.\nHowever, synthetic resin gears have drawbacks which metal gears do not, such as sink marks (i.e., defective shallow depressions) resulting from material shrinkage after a thick gear is molded. It is therefore necessary for synthetic resin gears to have a minimum thickness in order to prevent such sink marks from forming, and at the same time have sufficient rigidity and strength.\nIn a conventional gear 100 shown in FIGS. 5 and 6, an annular rib 106 is disposed between a boss 102 formed at the center of the gear 100 and a rim 104. Ribs 108 are interposed between and connect the boss 102 and the annular rib 106, and ribs 110 are interposed between and connect the annular rib 106 and the rim 104. A disc-shaped member 112 is disposed at a width-direction center of the rim 104 and connects (extends between) the rim 104 and the boss 102. Teeth are disposed along an outer periphery of the rim 104, with bottom lands being formed between respectively adjacent teeth and top lands 114 being formed at crests of the teeth. A radius R is equal to the distance from the center of the boss 102 to an imaginary circle formed by joining the top lands 114 (hereinafter, “top land circle”).\nAs mentioned above, when the gear 100 is made of synthetic resin, there arises a problem in that portions of top lands 114 corresponding to positions at which the ribs 110 connect to an inner peripheral surface of the rim 104 may be pulled radially inwards due to material shrinkage of the ribs 110 and the like. As a result, radii R drawn to top lands 114 corresponding to the positions at which the ribs 110 connect to the rim 104 become smaller than radii R drawn to top lands 114 not corresponding to positions at which the ribs 110 connect to the rim 104 (i.e., the gear 100 becomes radially nonuniform), whereby the top land circle becomes distorted.\nWhen a synthetic resin gear having a distorted top land circle is used, angle speed of a member (e.g., another gear or a timing belt) driven by the gear varies despite the angle speed of the gear being constant, which results in uneven rotation of the driven member. Consequently, uneven rotation of the driven member can adversely affect the drive in which the gear is employed."}
{"text": "The invention relates to novel phenyl-substituted cyclic ketoenols, to a plurality of processes and intermediates for their preparation and to their use as pesticides.\nIt has akready been discosed that certain phenyl-substituted cyc ketoenols are active as insectidces, acaricides and/or herbicides.\nAlso known are 1H-arylpyrrolidine-dione derivatives (BP-A-456 063, EP-A-521 334, EP-A-596 298, EP-A-613 884, EP-A-613 885, DE 44 40 594, WO 94/01 997, WO 95/01 358, WO 95126 954, WO 95/20 572, EP-A-0 668 267, WO 96/15 395, WO 96/35 664, WO 97/01 535 and WO 97/02 243).\nIt is known that certain substituted xcex943-dihydrofuran-2-one derivatives have herbicidal properties (cf. DE-A-4 014 420). The synthesis of the tetronic acid derivatives used as staring materials (such as, for example, 3-(2-methyl-phenyl)-4-hydroxy-5-(4-fluorophenyl)-xcex943-dihydrofuran-2-one) is likewise described in DE-A-4 014 420. Compounds of similar struftre without details of an insecticidal and/or acaricidal activity are known from the publication Campbell et al., J. Chem. Soc., Perkin Trans. 1, 1985, (S) 1567-76. Furthermore, 3-aryl-xcex943-dihydrofuranone derivatives having herbicidal, acaricidal and insecticidal properties are disclosed in EP-A-528 156, EP-A 0 647 637, WO 95/26345, WO 96/20 196, WO 96/25 395, WO 96/35 664, WO 97/01 535 and WO 97/02 243.\nCertain phenytl-pyrone derivatives unsubstituted in the phenyl ring have already been dislosed (cf. AM. Chirazi, T. Kappe and E. Ziegler, Arch. Phatin. 309, 558 (1976) and K. -H. Boltze and K. Heidenbluth, Chem. Ber. 91, 2849), a possible utility for these compounds as pesticides not being indicated. Phenyl-pyrone dertves substituted in the phenyl ring and having herbicidal, acaricidal and insecticidal propeties are descnrbed in EP-A-588 137, WO 96/25 395, WO 96/35 664 WO 97/01 535 and WO 97/02 243.\nHowever, the acaricidal and insecticidal activity and/or spectrum of action, and/or the toleration of the known compounds by plants, in parficular by crops, is not always satisfactory.\nThe invention, accordingly, provides compounds of the formula (I) \nin which\nX represents halogen, alkyl, alkenyl, alkynyl, alkoxy, benzyloxy, halogenoalkyl, halogenoalkoxy, cyano or nitro,\nZ represents hydrogen, amino, halogen, alkyl, alkoxy, halogenoalkyl, halogenoalkoxy, hydroxyl, cyano, nitro or respectively optionally substituted phenoxy, phenylthio, 5- or 6-membered hetaryloxy, 5- or 6-membered hetarylthio, phenylalkyloxy or phenylalkylthio and\nHet represents one of the groups \nxe2x80x83in which\nA represents a respectively optionally substituted radical from the group consisting of alkyl, alkenyl, alkoxyalkyl, polyalkoxyalkyl and alkylthioalkyl, represents respectively saturated or unsaturated and optionally substituted cycloalkyl or heterocyclyl or represents respectively optionally halogen-, alkyl-, halogenoalkyl-, alkoxy-, hatogenoalkoxy-, cyano- or nitro-substituted aryl, arylalkyl or hetaryl,\nB represents alkyl or alkoxyalkyl or\nA and B together with the carbon atom that they are attached to represent a saturated or unsaturated, optionally substituted carbocycle or heterocycle,\nD represents hydrogen or represents an optionally substituted radical from the group consisting of alkyl, alkenyl, alkynyl, alkoxyalkyl, polyalkoxyalkyl, alkylthioalkyl, saturated or unsaturated cycloalkyl, saturated or unsaturated heterocyclyl, arylalkyl, aryl, hetarylalkyl or hetaryl or\nA and D together with the atoms that they are attached to represent a respectively optionally substituted carbocycle or heterocycle,\nG represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur,\nM represents oxygen or sulphur,\nR1 represents respectively optionally halogen-substituted alkyl, alkenyl, alkoxyalkyl, alkylthioalkyl or polyalkoxyalkyl or represents respectively optionally halogen-, alkyl- or alkoxy-substituted cycloalkyl or heterocyclyl or represents respectively optionally substituted phenyl, phenylalkyl, hetaryl, phenoxyalkyl or hetarytoxyalkyl,\nR2 represents respectively optionally halogen-substituted alkyl, alkenyl, alkoxyalkyl or polyalkoxyalkyl or represents respectively optionally substituted cycloalkyl, phenyl or benzyl,\nR3, R4 and R5 independently of one another each represent respectively optionally halogen-substituted alkyl, alkoxy, alkylamino, dialkylamino, alkylthio, alkenylthio or cycloalkylthio or represent respectively optionally substituted phenyl, benzyl, phenoxy or phenylthio,\nR6 and R7 independently of one another each represent hydrogen, represent respectively optionally halogen-substituted alkyl, cycloalkyl, alkenyl, alkoxy, alkoxyalkyl, represent respectively optionally substituted phenyl or benzyl, or form together with the nitrogen atom that they are attached to an optionally oxygen- or sulphur-containing, optionally substituted cycle.\nThe compounds of the formula (I) can also be present, depending on the nature of the substituents, as geometric and/or optical isomers and isomer mixtures of differing composition which, if appropriate, can be separated in a customary manner. Both the pure isomers and the isomer mixtures, their preparation and use, and compositions comprising them are part of the subject matter of the present invention. In the following, for simplicity, however, compounds of the formula (I) are always referred to, although both pure compounds and, if appropriate, mixtures having different proportions of isomer compounds are intended.\nIncluding the meanings (1) to (4) of the group Het, the following principal structures (I-1) to (I-4) result: \nin which\nA, B, D, G, X and Z are each as defined above.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-1-a) to (I-1-g) result if Het represents the group (1) \nin which\nA, B, E, L, M, X, Z, R1, R2, R3, R4, R5, R6 and R7 are each as defined above.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-2-a) to (I-2-g) result if Het represents the group (2) \nin which\nA, B, E, L, M, X, Z, R1, R2, R3, R4, R5, R6 and R7 are each as defined above.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-3-a) to (I-3-g) result if Het represents the group (3) \nin which\nA, B, E, L, M, X, Z, R1, R2 R3, R4, R5, R6 and R7 are each as defined above.\nDepending on the position of the substituent G, the compounds of the formula (I-4) can be present in the two isomeric forms of formulae (I-4)a and (I-4)b\nwhich is intended to be expressed by the dashed line in the formula (I-4).\nThe compounds of the formulae (I-4)a and (I-4)b can be present both as mixtures and in the form of their pure isomers. Mixtures of the compounds of the formulae (I4)a and (I-4)b can, if desired, be separated by physical methods in a manner known per se, for example by chromatographic methods.\nFor better clarity, in the following in each case only one of the possible isomers is shown. This does not exclude the possibility that the compounds can optionally be present in the form of the isomer mixtures or in the other respective isomer form.\nIncluding the various meanings (a), (b), (c), (d), (e), (f) and (g) of the group G, the following principal structures (I-4-a) to (I-4-g) result if Het represents the group (4) \nin which\nA, D, E, L, M, X, Z, R1, R2, R3, R4, R5, R6 and R7 are each as defined above.\nFurthermore, it has been found that the novel compounds of the formula (I) are 10 obtained by one of the processes described below:\n(A) Compounds of the formula (I-1-a) \nin which\nA, B, X and Z are each as defined above, are obtained by the intramolecular condensation of compounds of the formula (II) \nxe2x80x83in which\nA, B, X and Z are each as defined above, and\nR8 represents alkyl (preferably C1-C6-alkyl), in the presence of a diluent and in the presence of a base.\n(B) Furthermore, it was found that compounds of the formula (I-2-a) \nin which\nA, B, X and Z are each as defined above,\nare obtained by the intramolecular condensation of compounds of the formula (III) \nxe2x80x83in which\nA, B, X, Z and R8 are each as defined above,\nin the presence of a diluent and in the presence of a base.\n(C) Furthermore, it was found that compounds of the formula (I-3-a) \nin which\nA, B, X and Z are each as defined above,\nare obtained by the intramolecular cyclization of compounds of the formula (IV) \nxe2x80x83in which\nA, B, X, Z and R8 are each as defined above and\nW represents hydrogen, halogen, alkyl (preferably C1-C6-alkyl) or alkoxy (preferably C1-C8-alkoxy), if appropriate in the presence of a diluent and in the presence of an acid.\n(D) Furthermore, it was found that compounds of the formula (I-4-a) \nin which\nA, D, X and Z are each as defined above,\nare obtained by reacting compounds of the formula (V) \nxe2x80x83in which\nA and D are each as defined above,\nor their silyl enol ethers of the formula (Va) \nxe2x80x83in which\nA and D are each as defined above and\nR8xe2x80x2 represents alkyl (preferably methyl),\nwith compounds of the formula (VI) \nxe2x80x83in which\nX and Z are each as defined above and\nHal represents halogen (preferably chlorine or bromine),\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid acceptor.\nFurthermore, it was found\n(E) that the compounds of the formulae (I-1-b) to (I-4-b) shown above in which A, B, D, R1, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above\nxcex1) with acid halides of the formula (VII) \nxe2x80x83in which\nR1 is as defined above and\nHal represents halogen (in particular chlorine or bromine), or\nxcex2) with carboxylic anhydrides of the formula (VIII)\nR1xe2x80x94COxe2x80x94Oxe2x80x94COxe2x80x94R1xe2x80x83xe2x80x83(VIII)\nxe2x80x83in which\nR1 is as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent;\n(F) that the compounds of the formulae (I-1-c) to (I-4-c) shown above in which A, B, D, R2, M, X and Z are each as defined above and L represents oxygen are obtained by reacting compounds of the formulae (I-1-a) to\n(I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith chloroformic esters or chloroformic thioesters of the formula (IX)\nR2xe2x80x94Mxe2x80x94COxe2x80x94Clxe2x80x83xe2x80x83(IX)\nxe2x80x83in which\nR2 and M are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent;\n(G) that compounds of the formulae (I-1-c) to (I-4-c) shown above in which A, B, D, R2, M, X and Z are each as defined above and L represents sulphur are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith chloromonothioformic esters or chlorodithioformic esters of the formula (X) \nxe2x80x83in which\nM and R2 are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-bindina agent,\n(H) that compounds of the formulae (I-1-d) to (I-4-d) shown above in which A, B, D, R3, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith sulphonyl chlorides of the formula (XII)\nR3xe2x80x94SO2xe2x80x94Clxe2x80x83xe2x80x83(XII)\nxe2x80x83in which\nR3 is as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent,\n(I) that compounds of the formulae (I-1-e) to (I-4-e) shown above in which A, B, D, L, R4, R5, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nwith phosphorus compounds of the formula (XIII) \nxe2x80x83in which\nL, R4 and R5 are each as defined above and\nHal represents halogen (in particular chlorine or bromine),\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent,\n(J) that compounds of the formulae (I-1-f) to (I-4-f) shown above in which A, B, D, E, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) in which A, B, D, X and Z are each as defined above, in each case\nwith metal compounds or amines of the formulae (XIV) or (XV) \nxe2x80x83in which\nMe represents a mono- or divalent metal (preferably an alkali metal or alkaline earth metal, such as lithium, sodium, potassium, magnesium or calcium),\nt represents the number 1 or 2 and\nR10, R11, R12 independently of one another each represent hydrogen or alkyl (preferably C1-C8-alkyl),\nif appropriate in the presence of a diluent;\n(K) that compounds of the formulae (I-1-g) to (I-4-g) shown above in which A, B, D, L, R6, R7, X and Z are each as defined above are obtained by reacting compounds of the formulae (I-1-a) to (I-4-a) shown above in which A, B, D, X and Z are each as defined above, in each case\nxcex1) with isocyanates or isothiocyanates of the formula (XVI)\nR6xe2x80x94Nxe2x95x90Cxe2x95x90Lxe2x80x83xe2x80x83(XVI)\nxe2x80x83in which\nR6 and L are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of a catalyst, or\nxcex2) with carbamoyl chlorides or thiocarbamoyl chlorides of the formula (XVII) \nxe2x80x83in which\nL, R6 and R7 are each as defined above,\nif appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nFurthermore, it has been found that the novel compounds of the formula (I) have a very good activity as pesticides, preferably as insecticides and acaricides, and that they additionally are very well tolerated by plants, in particular by crops.\nFormula (I) provides a general definition of the compounds according to the invention. Preferred substituents or ranges of the radicals shown in the formulae mentioned hereinabove and hereinbelow are illustrated below:\nX preferably represents halogen, C1-C6-alkyl, C2-C6-alkenyl, C2-C6-alkynyl, C1-C6-alkoxy, benzyloxy, C1-C4-halogenoalkyl, C1-C4-halogenoalkoxy, cyano or nitro.\nZ preferably represents hydrogen, amino, halogen, C1-C6-alkyl, C1-C6-alkoxy, C1-C4-halogenoalkyl, C1-C4-halogenoalkoxy, hydroxyl, cyano, nitro or respectively optionally halogen-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenoxy, phenylthio, thiazolyloxy, pyridinyloxy, pyrimidinyloxy, pyrazolyloxy, phenyl-C1-C4-alkyloxy or phenyl-C1-C7-alkylthio.\nHet preferably represents one of the groups \nA preferably represents respectively optionally halogen-substituted C1-C12-alkyl, C2-C8-alkenyl, C1-C10-alkoxy-C1-C8-alkyl, poly-C1-C8-alkoxy-C1-C8-alkyl or C1-C10-alkylthio-C1-C6-alkyl, preferably represents optionally halogen-, C1-C6-alkyl- or C1-C6-alkoxy-substituted C3-C8-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or preferably represents respectively optionally halogen-, C1-C6-alkyl-, C1-C6-halogenoalkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkoxy-, cyano- or nitro-substituted phenyl, naphthyl, phenyl-C1-C6-alkyl, naphthyl-C1-C6-alkyl or hetaryl having 5 or 6 ring atoms and one to three hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of furanyl, pyridyl, imidazolyl, triazolyl, pyrazolyl, indolyl, thiazolyl and thienyl).\nB preferably represents C1-C12-alkyl or C1-C8-alkoxy-C1-C6-alkyl or\nA, B and the carbon atom that they are attached to preferably represent C3-C10-cycloalkyl or C1-C10-cycloalkenyl where in each case one methylene group is optionally replaced by oxygen or sulphur and which are optionally substituted by C1-C8-alkyl, C3-C10-cycloalkyl, C1-C8-halogenoalkyl, C1-C8-alkoxy, C1-C8-alkylthio, halogen or phenyl or\nA, B and the carbon atom that they are attached to preferably represent C5-C6-cycloalkyl which is substituted by an alkylenediyl group optionally containing one or two not directly adjacent oxygen and/or sulphur atoms, or by an alkylenedioxy group or an alkylenedithioyl group forming a further five- to eight-membered ring with the carbon atom that it is attached to, or\nA, B and the carbon atom that they are attached to preferably represent C3-C8-cycloalkyl or C5-C8-cycloalkenyl in which two carbon atoms are linked to each other by respectively optionally C1-C6-alkyl-, C1-C6-alkoxy- or halogen-substituted C3-C6-alkanediyl, C1-C6-alkenediyl or C4-C6-alkanedienediyl, one methylene group in each case being optionally replaced by oxygen or sulphur.\nD preferably represents hydrogen, preferably represents respectively optionally halogen-substituted C1-C12-alkyl, C3-C8-alkenyl, C3-C8-alkynyl, C1-C10-alkoxy-C2-C8-alkyl, poly-C1-C8-alkoxy-C2-C8-alkyl or C1-C10-alkylthio-C2-C8-alkyl, preferably represents optionally halogen-, C1-C4-alkyl-, C1-C4-alkoxy- or C1-C4-halogenoalkyl-substituted C3-C8-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or preferably represents respectively optionally halogen-, C1-C6 alkyl-, C1-C6-halogenoalkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkoxy-, cyano- or nitro-substituted phenyl, hetaryl having 5 or 6 ring atoms and one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of furanyl, imidazolyl, pyridyl, thiazolyl, pyrazolyl, pyrimidyl, pyrrolyl, thienyl and triazolyl), phenyl-C1-C6-alkyl or hetaryl-C1-C6-alkyl having 5 or 6 ring atoms and one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of furanyl-, imidazolyl-, pyridyl-, thiazolyl-, pyrazolyl-, pyrimidyl-, pyrrolyl-, thienyl- and triazolyl-C1-C6-alkyl) or\nA and D together preferably represent a C3-C6-alkanediyl, C3-C6-alkenediyl or C4-C6-alkadienediyl group in which respectively optionally one methylene group is replaced by oxygen or sulphur and which are respectively optionally substituted by halogen or respectively optionally halogen-substituted C1-C10-alkyl, C1-C6-alkoxy, C1-C6-alkylthio, C3-C7-cycloalkyl, phenyl or benzyloxy or by a further C3-C6-alkanediyl, C3-C6-alkenediyl or C4-C6-alkadienediyl group forming a fused ring, in which optionally respectively one methylene group is replaced by oxygen or sulphur and which are optionally substituted by C1-C6-alkyl, or\nA and D together represent a C3-C6-alkanediyl or C3-C6-alkenediyl group containing in each case optionally one of the following groups \nG preferably represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur and\nM represents oxygen or sulphur.\nR1 preferably represents respectively optionally halogen-substituted C1-C20-alkyl, C2-C20-alkenyl, C1-C8-alkoxy-C1-C8-alkyl, C1-C8-alkylthio-C1-C8-alkyl or poly-C1-C8-alkoxy-C1-C8-alkyl or preferably represents optionally halogen-, C1-C6-alkyl- or C1-C6-alkoxy-substituted C3-C8-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur,\npreferably represents optionally halogen-, cyano-, nitro-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkyl-, C1-C6-halogenoalkoxy-, C1-C6-alkylthio- or C1-C6-alkylsulphonyl-substituted phenyl,\npreferably represents optionally halogen-, nitro-, cyano-, C1-C6-alkyl, C1-C6-alkoxy-, C1-C6-halogenoalkyl- or C1-C6-halogenoalkoxy-substituted phenyl-C1-C6-alkyl,\npreferably represents optionally halogen- or C1-C6-alkyl-substituted 5- or 6-membered hetaryl having one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of pyrazolyl, thiazolyl, pyridyl, pyrimidyl, furanyl and thienyl), preferably represents optionally halogen- or C1-C6-alkyl-substituted phenoxy-C1-C6-alkyl or\npreferably represents optionally halogen-, amino- or C1-C6-alkyl-substituted 5- or 6-membered hetaryloxy-C1-C6-alkyl having one or two hetero atoms from the group consisting of oxygen, sulphur and nitrogen (in particular from the group consisting of pyridyloxy-C1-C6-alkyl, pyrimidyloxy-C1-C6-alkyl and thiazolyloxy-Ci-C6-alkyl).\nR2 preferably represents respectively optionally halogen-substituted C1-C20-alkyl, C2-C20-alkenyl, C1-C8-alkoxy-C2-C8-alkyl or poly-C1-C8-alkoxy-C2-C8-alkyl,\npreferably represents optionally halogen-, C1-C6-alkyl- or C1-C6-alkoxy-substituted C3-C8-cycloalkyl or\npreferably represents respectively optionally halogen-, cyano-, nitro-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C6-halogenoalkyl- or C1-C6-halogenoalkoxy-substituted phenyl or benzyl.\nR3 preferably represents optionally halogen-substituted C1-C8-alkyl or respectively optionally halogen-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, cyano- or nitro-substituted phenyl or benzyl.\nR4 and R5 independently of one another each preferably represent respectively optionally halogen-substituted C1-C8-alkyl, C1-C8-alkoxy, C1-C8-alkylamino, di-(C1-C8-alkyl)-amino, C1-C8-alkylthio or C3-C8-alkenylthio or preferably represent respectively optionally halogen-, nitro-, cyano-, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, C1-C4-alkylthio-, C1-C4-halogenoalkylthio-, C1-C4-alkyl- or C1-C4-halogenoalkyl-substituted phenyl, phenoxy or phenylthio.\nR6 and R7 independently of one another each preferably represent hydrogen, preferably represent respectively optionally halogen-substituted C1-C8-alkyl, C3-C8-cycloalkyl, C1-C8-alkoxy, C3-C8-alkenyl or C1-C8-alkoxy-C2-C8-alkyl, preferably represent respectively optionally halogen-, C1-C8-alkyl-, C1-C8-halogenoalkyl- or C1-C8-alkoxy-substituted phenyl or benzyl or together preferably represent an optionally C1-C6-alkyl-substituted C3-C6-alkylene radical in which optionally one methylene group is replaced by oxygen or sulphur.\nR13 preferably represents hydrogen or respectively optionally halogen-substituted C1-C8-alkyl or C1-C8-alkoxy, preferably represents optionally halogen-, C1-C4-alkyl- or C1-C4-alkoxy-substituted C3-C8-cycloalkyl in which optionally one methylene group is replaced by oxygen or sulphur, or preferably represents respectively optionally halogen-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenyl, phenyl-C1-C4-alkyl or phenyl-C1-C4-alkoxy.\nR14 preferably represents hydrogen or C1-C8-alkyl or\nR13 and R14 together preferably represent C4-C6-alkanediyl.\nR15 and R16 are identical or different and each preferably represent C1-C6-alkyl or\nR15 and R16 together preferably represent a C2-C4-alkanediyl radical which is optionally substituted by C1-C6-alkyl or by optionally halogen-, C1-C4-alkyl-, C1-C4-halogenoalkyl-, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenyl.\nR17 and R18 independently of one another each preferably represent hydrogen, preferably represent optionally halogen-substituted C1-C8-alkyl or preferably represent optionally halogen-, C1-C6-alkyl-, C1-C6-alkoxy-, C1-C4-halogenoalkyl-, C1-C4-halogenoalkoxy-, nitro- or cyano-substituted phenyl or\nR17 and R18 together with the carbon atom that they are attached to preferably represent optionally C1-C4-alkyl-substituted C5-C7-cycloalkyl in which optionally one methylene group is replaced by oxygen or sulphur.\nR19 and R20 independently of one another each preferably represent C1-C10-alkyl, C1-C10-alkenyl, C1-C10-alkoxy, C1-C10-alkylamino, C3-C10-alkenylamino, di-(C1-C10-alkyl)-amino or di-(C3-C10-alkenyl)-amino.\nX particularly preferably represents fluorine, chlorine, bromine, C1-C4-alkyl, C1-C4-alkoxy, benzyloxy, C1-C2-halogenoalkyl, C1-C2-halogenoalkoxy, cyano or nitro.\nZ particularly preferably represents hydrogen, amino, fluorine, chlorine, bromine, C1-C4-alkyl, C1-C4-alkoxy, C1-C2-halogenoalkyl, C1-C2-halogenoalkoxy, hydroxyl, cyano, nitro or respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C2-halogenoalkyl-, C1-C2-halogenoalkoxy-, nitro- or cyano-substituted phenoxy or benzyloxy.\nHet particularly preferably represents one of the groups \nA particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C10-alkyl, C2-C6-alkenyl, C1-C8-alkoxy-C1-C6-alkyl, poly-C1-C6-alkoxy-C1-C6-alkyl or C1-C8-alkylthio-C1-C6-alkyl or particularly preferably represents optionally fluorine-, chlorine-, C1-C4-alkyl- or C1-C4-alkoxy-substituted C3-C7-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-halogenoalkyl-, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, cyano, or nitro-substituted phenyl, furanyl, pyridyl, imidazolyl, triazolyl, pyrazolyl, indolyl, thiazolyl, thienyl or phenyl-C1-C4-alkyl.\nB particularly preferably represents C1-C10-alkyl or C1-C6-alkoxy-C1-C4-alkyl or\nA, B and the carbon atom that they are attached to particularly preferably represent C3-C8-cycloalkyl or C5-C8-cycloalkenyl where one methylene group in each case is optionally replaced by oxygen or sulphur and which are optionally substituted by C1-C6-alkyl, C3-C8-cycloalkyl, C1-C3-haloogenoalkyl, C1-C6-alkoxy, C1-C6-alkylthio, fluorine, chlorine or phenyl or\nA, B and the carbon atom that they are attached to particularly preferably represent C1-C6-cycloalkyl which is substituted by an alkylenediyl group optionally containing one or two not directly adjacent oxygen or sulphur atoms or by an alkylenedioxy group or by an alkylenedithiol group which forms together with the carbon atom that it is attached to a further five- to seven-membered ring or\nA, B and the carbon atom that they are attached to particularly preferably represent C3-C6-cycloalkyl or C5-C6-cycloalkenyl in which two carbon atoms are linked to each other by respectively optionally C1-C4-alkyl-, C1-C4-alkoxy-, fluorine-, chlorine- or bromine-substituted C3-C5-alkanediyl or C3-C5-alkenediyl, one methylene group in each case being optionally replaced by oxygen or sulphur, or are linked to each other by butadienediyl.\nD particularly preferably represents hydrogen, particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl or C1-C6-alkoxy, particularly preferably represents optionally fluorine-, chlorine-, C1-C4-alkyl-, C1-C4-alkoxy- or C1-C2-halogenoalkyl-substituted C3-C7-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-halogenoalkyi, C1-C4-alkoxy-, C1-C4-halogenoalkoxy-, cyano- or nitro-substituted phenyl, furanyl, imidazolyl, pyridyl, thiazolyl, pyrazolyl, pyrimidyl, pyrrolyl, thienyl, triazolyl or phenyl-C1-C4-alkyl or\nA and D together particularly preferably represent a C3-C5-alkanediyl or C1-C5-alkenediyl group in which in each case one methylene group is optionally replaced by oxygen or sulphur and which are optionally substituted by fluorine, chlorine or by respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C1-C4-alkoxy, C1-C4-alkylthio, C3-C6-cycloalkyl, phenyl or benzyloxy or\nin which in each case optionally one of the following groups is contained: \nG particularly preferably represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur and\nM represents oxygen or sulphur.\nR1 particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C 16-alkyl, C2-C16-alkenyl, C1-C6-alkoxy-C1-C6-alkyl, C1-C6-alkylthio-C1-C6-alkyl or poly-C1-C6-alkoxy-C1-C6-alkyl or particularly preferably represents optionally fluorine-, chlorine-, C1-C5-alkyl- or C1-C5-alkoxy-substituted C3-C7-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur,\nparticularly preferably represents optionally fluorine-, chlorine-, bromine-, cyano-, nitro-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C3-halogenoalkyl-, C1-C3-halogenoalkoxy-, C1-C4-alkylthio- or C1-C4-alkylsulphonyl-substituted phenyl,\nparticularly preferably represents optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C3-halogenoalkyl- or C1-C3-halogeno-alkoxy-substituted phenyl-C1-C4-alkyl,\nparticularly preferably represents respectively optionally fluorine-, chlorine-, bromine- or C1-C4-alkyl-substituted pyrazolyl, thiazolyl, pyridyl, pyrimidyl, furanyl or thienyl,\nparticularly preferably represents optionally fluorine-, chlorine-, bromine- or C1-C4-alkyl-substituted phenoxy-C1-C5-alkyl or\nparticularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, amino- or C1-C4-alkyl-substituted pyridyloxy-C1-C5-alkyl, pyrimidyloxy-C1-C5-alkyl or thiazolyloxy-C1-C5-alkyl.\nR2 particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C16-alkyl, C2-C16-alkenyl, C1-C6-alkoxy-C2-C6-alkyl or poly-C1-C6-alkoxy-C2-C6-alkyl,\nparticularly preferably represents optionally fluorine-, chlorine-, C1-C4-alkyl- or C1-C4-alkoxy-substituted C3-C7-cycloalkyl or\nparticularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, cyano-, nitro-, C1-C4-alkyl-, C1-C3-alkoxy-, C1-C3-halogenoalkyl- or C1-C3-halogenoalkoxy-substituted phenyl or benzyl.\nR3 particularly preferably represents optionally fluorine- or chlorine-substituted C1-C6-alkyl or respectively optionally fluorine-, chlorine-, bromine-, C1-C4-alkyl-, C1-C4-alkoxy-, C1-C2-halogenoalkoxy-, C1-C2-halogenoalkyl-, cyano-, or nitro-substituted phenyl or benzyl.\nR4 and R5 independently of one another each particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C1-C6-alkoxy, C1-C6-alkylamino, di-(C1-C6-alkyl)-amino, C1-C6-alkylthio or C3-C4-alkenylthio or particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, nitro-, cyano-, C1-C3-alkoxy-, C1-C3-haloge-noalkoxy-, C1-C3-alkylthio-, C1-C3-halogenoalkylthio-, C1-C3-alkyl- or C1-C3-halogenoalkyl-substituted phenyl, phenoxy or phenylthio.\nR6 and R7 independently of one another each particularly preferably represent hydrogen, particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C3-C6-cycloalkyl, C1-C6-alkoxy, C3-C6-alkenyl or C1-C6-alkoxy-C2-C6-alkyl, particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, C1-C5-halogenoalkyl-, C1-C5-alkyl- or C1-C5-alkoxy-substituted phenyl or benzyl, or together particularly preferably represent an optionally C1-C4-alkyl-substituted C3-C6-alkylene radical in which optionally one methylene group is replaced by oxygen or sulphur.\nR13 particularly preferably represents hydrogen or respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl or C1-C6-alkoxy, particularly preferably represents optionally fluorine-, C1-C2-alkyl- or C1-C2-alkoxy-substituted C3-C7-cycloalkyl in which optionally one methylene group is replaced by oxygen or sulphur, or particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, C1-C5-alkyl-, C1-C5-alkoxy-, C1-C2-halogenoalkyl-, C1-C2-halogenoalkoxy-, nitro- or cyano-substituted phenyl, phenyl-C1-C3-alkyl or phenyl-C1-C2-alkyloxy.\nR14 particularly preferably represents hydrogen or C1-C6-alkyl or\nR13 and R14 together particularly preferably represent C4-C6-alkanediyl.\nR15 and R16 are identical or different and each particularly preferably represent C1-C4-alkyl or\nR15 and R16 together particularly preferably represent a C2-C3-alkanediyl radical which is optionally substituted by C1-C4-alkyl or by optionally fluorine-, chlorine-, bromine-, C1-C2-alkyl-, C1-C2-halogenoalkyl-, C1-C2-alkoxy-, C1-C2-halogenoalkoxy-, nitro- or cyano-substituted phenyl.\nX very particularly preferably represents fluorine, chlorine, bromine, methyl, ethyl, n-propyl, iso-propyl, n-butyl, iso-butyl, tert-butyl, methoxy, ethoxy, n-propoxy, iso-propoxy, trifluoromethyl, trifluoromethoxy, difluoromethoxy, cyano or nitro,\nZ very particularly preferably represents hydrogen, amino, fluorine, chlorine, bromine, methyl, ethyl, n-propyl, iso-propyl, n-butyl, iso-butyl, tert-butyl, methoxy, ethoxy, n-propoxy, iso-propoxy, trifluoromethyl, trifluoromethoxy, difluoromethoxy, cyano or nitro.\nHet very particularly preferably represents one of the groups \nA very particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C8-alkyl, C2-C4-alkenyl, C1-C6-alkoxy-C1-C4-alkyl, poly-C1-C4-alkoxy-C1-C4-alkyl or C1-C6-alkylthio-C1-C4-alkyl, or very particularly preferably represents optionally fluorine-, chlorine-, methyl- or methoxy-substituted C3-C6-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or very particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, n-propyl-, iso-propyl-, methoxy-, ethoxy-, trifluoromethyl-, trifluoromethoxy-, cyano- or nitro-substituted phenyl, pyridyl or benzyl.\nB very particularly preferably represents C1-C8-alkyl or C1-C4-alkoxy-C1-C2-alkyl or\nA, B and the carbon atom that they are attached to very particularly preferably represent C3-C8-cycloalkyl or C5-C8-cycloalkenyl in which in each case optionally one methylene group is replaced by oxygen or sulphur and which are optionally substituted by methyl, ethyl, n-propyl, iso-propyl, butyl, iso-butyl, sec-butyl, tert-butyl, cyclohexyl, trifluoromethyl, methoxy, ethoxy, n-propoxy, iso-propoxy, butoxy, iso-butoxy, sec-butoxy, tert-butoxy, methylthio, ethylthio, fluorine, chlorine or phenyl or\nA, B and the carbon atom that they are attached to very particularly preferably represent C3-C6-cycloalkyl or C5-C6-cycloalkenyl in which two carbon atoms are linked together by C3-C4-alkanediyl or C3-C4-alkenediyl in which in each case optionally one methylene group is replaced by oxygen or sulphur, or are linked together by butadienediyl.\nD very particularly preferably represents hydrogen, very particularly preferably represents respectively optionally fluorine-substituted C1-C4-alkyl or C1-C4-alkoxy or C3-C6-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur, or very particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, n-propyl-, iso-propyl-, methoxy-, ethoxy-, trifluoromethyl-, trifluoromethoxy-, cyano- or nitro- substituted phenyl, furanyl, pyridyl, thienyl or benzyl, or\nA and D together very particularly preferably represent a C3-C5-alkanediyl or C3-C5-alkenediyl group in which in each case optionally one methylene group is replaced by oxygen or sulphur and which are optionally substituted by fluorine, chlorine or by respectively optionally fluorine- or chlorine-substituted C1-C6-alkyl, C1-C4-alkoxy, C1-C4-alkylthio, C3-C6-cycloalkyl, phenyl or benzyloxy.\nG very particularly preferably represents hydrogen (a) or represents one of the groups \nxe2x80x83in which\nE represents a metal ion equivalent or an ammonium ion,\nL represents oxygen or sulphur and\nM represents oxygen or sulphur.\nR1 very particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C14-alkyl, C2-C14-alkenyl, C1-C4-alkoxy-C1-C6-alkyl, C1-C4-alkylthio-C1-C6-alkyl, poly-C1-C4-alkoxy-C1-C4-alkyl or very particularly preferably represents optionally fluorine-, chlorine-, methyl-, ethyl-, n-propyl-, i-propyl-, n-butyl-, i-butyl-, tert-butyl-, methoxy-, ethoxy-, n-propoxy- or iso-propoxy-substituted C3-C6-cycloalkyl in which optionally one or two not directly adjacent methylene groups are replaced by oxygen and/or sulphur,\nvery particularly preferably represents optionally fluorine-, chlorine-, bromine-, cyano-, nitro-, methyl-, ethyl-, n-propyl-, i-propyl-, methoxy-, ethoxy-, trifluoromethyl-, trifluoromethoxy-, methylthio-, ethylthio-, methylsulphonyl- or ethylsulphonyl-substituted phenyl,\nvery particularly preferably represents optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, n-propyl-, i-propyl-, methoxy-, ethoxy-, trifluoromethyl- or trifluoromethoxy-substituted benzyl,\nvery particularly preferably represents respectively optionally fluorine-, chlorine-, bromine-, methyl- or ethyl-substituted furanyl, thienyl or pyridyl,\nvery particularly preferably represents optionally fluorine-, chlorine-, methyl- or ethyl-substituted phenoxy-C1-C4-alkyl or\nvery particularly preferably represents respectively optionally fluorine-, chlorine-, amino-, methyl- or ethyl-substituted pyridyloxy-C1-C4-alkyl, pyrimidyloxy-C1-C4-alkyl or thiazolyloxy-C1-C4-alkyl.\nR2 very particularly preferably represents respectively optionally fluorine- or chlorine-substituted C1-C14-alkyl, C2-C14-alkenyl, C1-C4-alkoxy-C2-C6-alkyl or poly-C1-C4-alkoxy-C2-C6-alkyl,\nvery particularly preferably represents optionally fluorine-, chlorine-, methyl-, ethyl-, n-propyl-, iso-propyl- or methoxy-substituted C3-C6-cycloalkyl,\nor very particularly preferably represents respectively optionally fluorine-, chlorine-, cyano-, nitro-, methyl-, ethyl-, n-propyl-, i-propyl-, methoxy-, ethoxy-, trifluoromethyl- or trifluoromethoxy-substituted phenyl or benzyl.\nR3 very particularly preferably represents optionally fluorine- or chlorine-substituted methyl, ethyl, propyl, iso-propyl, n-butyl, tert-butyl or respectively optionally fluorine-, chlorine-, bromine-, methyl-, ethyl-, iso-propyl-, tert-butyl-, methoxy-, ethoxy-, iso-propoxy-, trifluoromethyl-, trifluoromethoxy-, cyano- or nitro-substituted phenyl or benzyl.\nR4 and R5 independently of one another each very particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C4-alkyl, C1-C4-alkoxy, C1-C4-alkylamino, di-(C1-C4-alkyl)-amino or C1-C4-alkylthio or very particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, nitro-, cyano-, methyl-, methoxy-, trifluoromethyl- or trifluoromethoxy-subsfituted phenyl, phenoxy or phenylthio.\nR6 and R7 independently of one another each very particularly preferably represent hydrogen, very particularly preferably represent respectively optionally fluorine- or chlorine-substituted C1-C4-alkyl, C3-C6-cycloalkyl, C1-C4-alkoxy, C3-C4-alkenyl or C1-C4-alkoxy-C2-C4-alkyl, very particularly preferably represent respectively optionally fluorine-, chlorine-, bromine-, methyl-, methoxy- or trifluoromethyl-substituted phenyl or benzyl, or together very particularly preferably represent an optionally methyl- or ethyl-substituted C5-C6-alkylene radical in which optionally one methylene group is replaced by oxygen or sulphur.\nParticular preference is given to compounds of the formula (I) in which Z does not represent hydrogen.\nPreference is also given to those compounds where D does not represent methyl.\nThe abovementioned definitions or illustrations of radicals mentioned generally or in preferred ranges can be combined with each other as desired, i.e. also between the respective ranges and preferred ranges. They apply correspondingly to the final products and to the precursors and intermediates.\nFor the purpose of the invention, preference is given to compounds of the formula (I) in which there exists a combination of the meanings mentioned above as preferred (preferably).\nFor the purpose of the invention, particular preference is given to compounds of the formula (I) in which there exists a combination of the meanings mentioned above as particularly preferred.\nFor the purpose of the invention, very particular preference is given to compounds of the formula (I) in which there exists a combination of the meanings mentioned above as very particularly preferred.\nSaturated or unsaturated hydrocarbon radicals such as alkyl or alkenyl can, as far as possible, in each case be straight-chain or branched, also in combination with hetero atoms, for example in alkoxy.\nOptionally substituted radicals can be mono- or polysubstituted, it being possible in the case of polysubstitution for the substituents to be identical or different.\nIn addition to the compounds mentioned in the Preparation Examples, the following compounds of the formula (I-1-a) may be specifically mentioned:\nIn addition to the compounds mentioned in the preparation examples, the following compounds of the formula (1-2-a) may be specifically mentioned:\nIn addition to the compounds mentioned in the preparation examples, the following compounds of the formula (I-3-a) may be specifically mentioned:\nIn addition to the compounds mentioned in the Preparation Examples, the following compounds of the formula (I-4-a) may be specifically mentioned:\nIf according to process (A) ethyl N-[(3-chloro-6-methyl)-phenylacetyl]-1-amino-4-ethyl-cyclohexane-carboxylate is used as starting material, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (B) ethyl O-[(2,5-dichloro)-phenylacetyl]-2-hydroxy-isobutyrate is used, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (C) ethyl 2-[(2-chloro-5-methyl)-phenyl]-4-(4-methoxy)-benzylmercapto-4-methyl-3-oxo-valerate is used, the course of the process according to the invention can be represented by the following reaction scheme: \nIf, for example, according to process (D) (chlorocarbonyl)-2-[(3-chloro-6-methyl)-phenyl]-ketene and 4-fluoro-propiophenone are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (Excex1) 3-[(2,5-dichloro)-phenyl]-5,5-dimethylpyrrolidine-2,4-dione and pivaloyl chloride are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (E) (variant xcex2) 3-[(2,5-dichloro)-phenyl]-4-hydroxy-5-methyl-5-phenyl-xcex943-dihydrofuran-2-one and acetic anhydride are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (F) 8-[(2-chloro-5-methyl)-phenyl]-5,5-pentamethylene-pyrrolidine-2,4-dione and ethoxyethyl chloroformate are used as starting materials, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (G) 3-[(2-bromo-5-methyl)-phenyl]-4-hydroxy-5-methyl-6-(3-pyridyl)-pyrone and methyl chloromonothioformate are used as starting materials, the course of the reaction can be represented in the following manner: \nIf according to process (H) 2-[(2,5-dimethyl)-phenyl]-5,5-[(3-methyl)-pentamethylene]pyrrolidine-2,4-dione and methanesulphonyl chloride are used as starting materials, the course of the reaction can be represented by the following reaction scheme: \nIf according to process (I) 2-[(2-chlroro-5methyl)-phenyl]-4-hydroxy-5-methyl-6-(2-pyridyl)-pyrone and 2,2,2-trifluoroethyl chloromethanethio-phosphonate are used as starting materials, the course of the reaction can be represented by the following reaction scheme: \nIf according to process (J) 3-[(2,5-dichloro)-phenyl]-5-cyclopropyl-5-methyl-pyrrolidine-2,4-dione and NaOH are used as components, the course of the process according to the invention can be represented by the following reaction scheme: \nIf according to process (K) (variant xcex1) 3-[(2-chloro-5-methyl)-phenyl]-4-hydroxy-5,5-tetramethylene-xcex943-dihydro-furan-2-one and ethyl isocyanate are used as starting materials, the course of the reaction can be represented by the following reaction scheme: \nIf according to process (K) (variant xcex2) 3-[(2-chloro-5-methyl)-phenyl]-5,5-dimethyl-pyrrolidine-2,4-dione and dimethylcarbamoyl chloride are used as starting materials, the course of the reaction can be represented by the following scheme: \nThe compounds of the formula (II) \nin which\n5 A, B, X, Z and R8 are each as defined above,\nrequired as starting materials in process (A) according to the invention are novel.\nThe acylamino acid esters of the formula (II) are obtained, for example, when amino acid derivatives of the formula (XVIII) \nin which\nA, B and R8 are as defined above,\nare acylated using substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX and Z are each as defined above and\nHal represents chlorine or bromine,\n(Chem. Reviews 52, 237-416 (1953); Bhattacharya, Indian J. Chem. 6, 341-5, 1968)\nor when acylamino acids of the formula (XX) \nxe2x80x83in which\nA, B, X and Z are each as defined above,\nare esterified (Chem. Ind. (London) 1568 (1968)).\nThe compounds of the formula (XX) \nin which\nA, B, X and Z are each as defined above,\nare novel.\nThe compounds of the formula (XX) are obtained when amino acids of the formula (XXI) \nin which\nA and B are each as defined above,\nare acylated using substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX and Z are each as defined above and\nHal represents chlorine or bromine,\naccording to Schotten-Baumann (Organikum, VEB Deutscher Verlag der Wissenschaften, Berlin 1977, p. 505).\nSome of the compounds of the formula (XIX) are novel. They can be prepared by known methods.\nThe compounds of the formula (XIX) are obtained, for example, by reacting substituted phenylacetic acids of the formula (XXII) \nin which\nX and Z are each as defined above,\nwith halogenating agents (for example thionyl chloride, thionyl bromide, oxalyl chloride, phosgene, phosphorus trichloride, phosphorus tribromide or phosphorus pentachloride), if appropriate in the presence of a diluent (for example optionally chlorinated aliphatic or aromatic hydrocarbons such as toluene or methylene chloride) at temperatures from xe2x88x9220xc2x0 C. to 150xc2x0 C., preferably from xe2x88x9210xc2x0 C. to 100xc2x0 C.\nThe compounds of the formula (XXII) are novel with the exception of 2,5-dichlorophenylacetic acid (CAS 5398798), 5-chloro-2-methoxyphenylacetic acid (CAS 7569-6-22), 2-chloro-5-methylphenylacetic acid (CAS 81682-39-5), 2,5-difluorophenylacetic acid (CAS 85068-27-5), 2-bromo-5-methylphenylacetic acid (BRN 3 249 577) and 2-chloro-5-trifluoromethylphenylacetic acid (CAS 22893-39-6), they can be prepared by methods known from the literature (Organikum, 15th edition, p. 533, VEB Deutscher Verlag der Wissenschaften, Berlin 1977). The compounds of the formula (XMI) are obtained, for example, by hydrolysing substituted phenylacetic acid esters of the formula (XXIII) \nin which\nX, Z and R8 are each as defined above,\nat temperatures between 0xc2x0 C. and 150xc2x0 C., preferably between 20xc2x0 C. and 100xc2x0 C., in the presence of an acid (for example an inorganic acid such as hydrochloric acid) or of a base (for example of an alkali metal hydroxide such as sodium hydroxide or potassium hydroxide) and, if appropriate, of a diluent (for example of an aqueous alcohol such as methanol or ethanol).\nThe compounds of the formula (XXIII) are novel with the exception of methyl 2,5-dichlorophenylacetate (CAS 96129-66-7) and methyl 5-chloro-2-methoxy-phenylacetate (CAS 26939-01-5), they can be prepared by methods known in principle.\nThe compounds of the formula (XXIII) are obtained, for example, by reacting substituted 1,1,1-trichloro-2-phenylethanes of the formula (XXIV) \nin which\nX and Z are each as defined above,\nfirst with alkoxides (for example alkali metal alkoxides such as sodium methoxide or sodium ethoxide) in the presence of a diluent (for example the alcohol derived from the alkoxide) at temperatures between 0xc2x0 C. and 150xc2x0 C., preferably between 20xc2x0 C. and 120xc2x0 C., and then reacting with an acid (preferably an inorganic acid, such as sulphuric acid) at temperatures between xe2x88x9220xc2x0 C. and 150xc2x0 C., preferably 0xc2x0 C. and 100xc2x0 C. (cf. DE 3 314 249).\nThe compounds of the formula (XXIV) are novel, they can be prepared by methods known in principle.\nThe compounds of the formula (XXIV) are obtained, for example, when anilines of the formula (XXV) \nin which\nX and Z are as defined above,\nare reacted with vinylidene chloride (CH2=CCl2) in the presence of an alkyl nitrite of the formula (XXVI)\nR21xe2x80x94ONOxe2x80x83xe2x80x83(XXVI)\nxe2x80x83in which\nR21 represents alkyl, preferably C1-C6-alkyl,\nin the presence of copper(II) chloride and, if appropriate, in the presence of a diluent (for example of an aliphatic nitrile such as acetonitrile) at a temperature of xe2x88x9220xc2x0 C. to 80xc2x0 C., preferably 0xc2x0 C. to 60xc2x0 C.\nThe compounds of the formulae (XXV) and (XXVI) are known compounds of organic chemistry. Copper(II) chloride and vinylidene chloride have long been known and are commercially available.\nSome of the compounds of the formulae (XVIII) and (XXI) are known and/or they can be prepared by known processes (see, for example, Compagnon, Miocque Ann. Chim. (Paris) [14]5, p. 11-22, 23-27 (1970)).\nThe substituted cyclic aminocarboxylic acids of the formula (XXIa), in which A and B form a ring, are in general obtainable by the Bucherer-Bergs synthesis or by the Strecker synthesis and are in each case obtained here in different isomeric forms. Thus, according to the conditions of the Bucherer-Bergs synthesis mainly the isomers (in the following designated as xcex2 for the sake of simplicity) in which the radicals R and the carboxyl group are equatorial are obtained, while according to the conditions of the Strecker synthesis mainly the isomers (in the following designated as xcex1 for the sake of simplicity) are obtained in which the amino group and the radicals R are equatorial. \n(L. Munday, J. Chem. Soc. 4372 (1961); J. T. Eward, C. Jitrangeri, Can. J. Chem. 53, 3339 (1975).\nFurthermore the starting materials of the formula (II) \nin which\nA, B, X, Z and R8 are each as defined above,\nused in the above process (A) can be prepared by reacting aminonitriles of the formula (XXVII) \nxe2x80x83in which\nA and B are each as defined above,\nwith substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX, Z and Hal are each as defined above,\nto give compounds of the formula (XXVIII) \nxe2x80x83in which\nA, B, X and Z are each as defined above,\nand then subjecting these to an acidic alcoholysis.\nThe compounds of the formula (XXVIII) are also novel.\nThe compounds of the formula (III) \nin which\nA, B, X, Z and R8 are each as defined above,\nrequired as starting materials in process (B) according to the invention are novel.\nThey can be prepared in a simple manner by methods known in principle.\nThe compounds of the formula (III) are obtained, for example, when 2-hydroxycarboxylic acid esters of the formula (XXIX) \nin which\nA, B and R8 are each as defined above,\nare acylated using substituted phenylacetyl halides of the formula (XIX) \nxe2x80x83in which\nX, Z and Hal are each as defined above\n(Chem. Reviews 52, 237-416 (1953)).\nThe compounds of the formula (IV) \nin which\nA, B, W, X, Z and R8 are each as defined above,\nrequired as starting materials in the above process (C) are novel.\nThey can be prepared by methods known in principle.\nThe compounds of the formula (IV) are obtained, for example, when substituted phenylacetic acid esters of the formula (XXIII) \nin which\nX, R8 and Z are each as defined above,\nare acylated using 2-benzylthio-carbonyl halides of the formula (XXX) \nin which\nA, B and W are each as defined above and\nHal represents halogen (in particular chlorine or bromine),\nin the presence of strong bases (see, for example, M. S. Chambers, E. J. Thomas, D. J. Williams, J. Chem. Soc. Chem. Commun., (1987), 1228).\nThe benzylthio-carbonyl halides of the formula (XXX) are known in some cases and/or can be prepared by known methods (J. Antibiotics (1983), 26, 1589).\nThe halogenocarbonylketenes of the formula (VI) in which Z does not represent hydrogen, which are required as starting materials in process (D), are novel. They can be prepared in a simple manner by methods known in principle (cf., for example, Org. Prep. Proced. Int., 7, (4), 155-158, 1975 and DE 1 945 703).\nThe compounds of the formula (VI) \nin which\nX and Z are each as defined above and\nHal represents chlorine or bromine,\nare obtained when substituted phenylmalonic acids of the formula (XXXI) \nxe2x80x83in which\nX and Z are each as defined above,\nare reacted with acid halides, for example thionyl chloride, phosphorus(V) chloride, phosphorus(III) chloride, oxalyl chloride, phosgene or thionyl bromide, if appropriate in the presence of catalysts, for example diethylformamide, methylstearylformamide or triphenylphosphine and, if appropriate, in the presence of bases, for example pyridine or triethylamine, at a temperature between xe2x88x9220xc2x0 C. and 200xc2x0 C., preferably between 0xc2x0 C. and 150xc2x0 C.\nThe substituted phenylmalonic acids of the formula (XXXI) in which Z does not represent hydrogen are novel. However, they may be prepared in a simple manner by known processes (cf., for example, Organikum, VEB Deutscher Verlag der Wissenschaften, Berlin 1977, p. 517 ff), for example by hydrolysis of substituted phenylmalonic esters of the formula (XXXI) \nin which\nX, Z and R8 are each as defined above.\nThe carbonyl compounds of the formula (V) or their silyl enole ethers of the formula (Va) \nin which\nA, D and R8 xe2x80x2are each as defined above,\nrequired as starting materials for process (D) according to the invention are compounds which are commercially available, generally known or accessible by known processes.\nThe malonic acid esters of the formula (XXXII) \nin which\nR8, X and Z are each as defined above, and Z is not hydrogen, are novel.\nThey can be prepared by generally known methods of organic chemistry (cf., for example, Tetrahedron Lett. 27, 2763 (1986) and Organikum VEB Deutscher Verlag der Wissenschaften, Berlin 1977, p.587 ff).\nThe acid halides of the formula (VII), carboxylic anhydrides of the formula (VIII), chloroformic acid esters or chloroformic acid thioesters of the formula (IX), chloromonothioformic acid esters or chlorodithioformic acid esters of the formula (X), sulphonyl chlorides of the formula (XII), phosphorus compounds of the formula (XIII) and metal hydroxides, metal alkoxides or amines of the formula (XIV) and (XV) and isocyanates of the formula (XVI) and carbamoyl chlorides of the formula (XVII) additionally required as starting materials for carrying out processes (F), (G), (H), (I), (J) and (K) according to the invention are generally known compounds of organic or inorganic chemistry.\nThe compounds of the formulae (V), (VII) to (XVII), (XVIII), (XXI), (XXII), (XXIX), (XXX) and (XXXI) are moreover disclosed in the patent applications cited at the outset and/or can be prepared by the methods given there.\nProcess (A) is characterized in that compounds of the formula (II) in which A, B, X, Z and R8 are each as defined above are subjected to an intramolecular condensation in the presence of a diluent and in the presence of a base.\nSuitable diluents for the process (A) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as toluene and xylene, furthermore ethers, such as dibutyl ether, tetrahydrofuran, dioxane, glycol dimethyl ether and diglycol dimethyl ether, additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethyl-formamide and n-methyl-pyrrolidone, and also alcohols such as methanol, ethanol, propanol, iso-propanol, butanol, iso-butanol and tert-butanol.\nSuitable bases (deprotonating agents) for carrying out process (A) according to the invention are all customary proton acceptors. Those preferably utilizable are alkali metal and alkaline earth metal oxides, hydroxides and carbonates, such as sodium hydroxide, potassium hydroxide, magnesium oxide, calcium oxide, sodium carbonate, potassium carbonate and calcium carbonate, each of which can also be employed in the presence of a phase-transfer catalyst, for example triethylbenzylammonium chloride, tetrabutylammonium bromide, Adogene 464 (=methyltrialkyl(C8-C10)ammonium chloride) or TDA 1 (=tris-(methoxyethoxyethyl)-amine. Alkali metals such as sodium or potassium can also be used. Furthermore, alkali metal and alkaline earth metal amides and hydrides, such as sodium amide, sodium hydride and calcium hydride, and additionally also alkali metal alkoxides, such as sodium methoxide, sodium ethoxide and potassium tert-butoxide can be employed.\nWhen carrying out process (A) according to the invention, the reaction temperature can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 150xc2x0 C.\nProcess (A) according to the invention is generally carried out under atmospheric pressure.\nWhen carrying out process (A) according to the invention, the reaction component of the formula (II) and the deprotonating base are generally employed in equimolar to approximately double equimolar amounts. However, it is also possible to use one component or the other in a relatively large excess (up to 3 mol).\nProcess (B) is characterized in that compounds of the formula (III) in which A, B, X, Z and R8 are each as defined above are condensed intramolecularly in the presence of a diluent and in the presence of a base.\nSuitable diluents for the process (B) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as toluene and xylene, furthermore ethers, such as dibutyl ether, tetrahydrofuran, dioxane, glycol dimethyl ether and diglycol dimethyl ether, and additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethylformamide and N-methyl-pyrrolidone. Alcohols such as methanol, ethanol, propanol, iso-propanol, butanol, iso-butanol and tert-butanol can also be used.\nSuitable bases (deprotonating agents) for carrying out the process (B) according to the invention are all customary proton acceptors. Those preferably utilizable are alkali metal and alkaline earth metal oxides, hydroxides and carbonates, such as sodium hydroxide, potassium hydroxide, magnesium oxide, calcium oxide, sodium carbonate, potassium carbonate and calcium carbonate, each of which can also be employed in the presence of phase-transfer catalysts, for example triethylbenzylammonium chloride, tetrabutylammonium bromide, Adogene 464 (=methyltrialkyl(C8-C10-ammonium chloride) or TDA 1 (=tris-(methoxyethoxyethyl)-amine). Alkali metals such as sodium or potassium can also be used. Suitable are also alkali metal and alkaline earth metal amides and hydrides, such as sodium amide, sodium hydride and calcium hydride, and also alkali metal alkoxides, such as sodium methoxide, sodium ethoxide and potassium tert-butoxide.\nWhen carrying out process (B) according to the invention, the reaction temperature can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 150xc2x0 C.\nProcess (B) according to the invention is generally carried out under atmospheric pressure.\nWhen carrying out the process (B) according to the invention, the reaction components of the formula (III) and the deprotonating bases are generally employed in approximately equimolar amounts. However, it is also possible to use one component or the other in a relatively large excess (up to 3 mol).\nProcess (C) is characterized in that compounds of the formula (IV) in which A, B, W, X, Z and R8 are each as defined above are cyclized intramolecularly in the presence of an acid and, if appropriate, in the presence of a diluent.\nSuitable diluents for the process (C) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as toluene and xylene, furthermore halogenated hydrocarbons, such as dichloromethane, chloroform, ethylene chloride, chlorobenzene, dichlorobenzene, and additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethylformamide and N-methyl-pyrrolidone. Alcohols such as methanol, ethanol, propanol, iso-propanol, butanol, iso-butanol and tert-butanol can also be used.\nThe acid employed can, if appropriate, also be used as a diluent.\nAcids which can be employed in process (C) according to the invention are all customary inorganic and organic acids, for example hydrohalic acids, sulphuric acid, alkyl-, aryl- and haloalkyl sulphonic acids; halogenated alkylcarboxylic acids, for example trifluoroacetic acid, are used in particular.\nWhen carrying out process (C) according to the invention, the reaction temperature can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 150xc2x0 C.\nProcess (C) according to the invention is generally carried out under atmospheric pressure.\nWhen carrying out process (C) according to the invention, the reaction components of the formulae (IV) and the acid are employed, for example, in equimolar amounts. However, it is, if appropriate, also possible to employ the acid in catalytic amounts.\nProcess (D) according to the invention is characterized in that carbonyl compounds of the formula (V) or their silyl enol ethers of the formula (Va) are reacted with ketene acid halides of the formula (VI), if appropriate in the presence of a diluent and if appropriate in the presence of an acid acceptor.\nSuitable diluents for the process (D) according to the invention are all organic solvents which are inert to the reaction participants. Those preferably utilizable are hydrocarbons, such as o-dichlorobenzene, tetraline, toluene and xylene, furthermore ethers, such as dibutyl ether, glycol dimethyl ether and diglycol dimethyl ether, and additionally polar solvents, such as dimethyl sulphoxide, sulpholane, dimethylformamide or N-methyl-pyrrolidone.\nAcid acceptors which can be used when carrying out process (D) according to the invention are all customary acid acceptors.\nThose preferably utilizable are tertiary amines, such as triethylamine, pyridine, diazabicyclooctane (DABCO), diazabicycloundecene (DBU), diazabicyclononene (DBN), Hunig base or N,N-dimethyl-aniline.\nWhen carrying out process (D) according to the invention, the reaction temperature can be varied within a relatively wide range. The reaction is expediently carried out at temperatures between 0xc2x0 C. and 250xc2x0 C., preferably between 50xc2x0 C. and 220xc2x0 C.\nProcess (D) according to the invention is preferably carried out under atmospheric pressure.\nWhen carrying out process (D) according to the invention, the reaction components of the formulae (V) and (VI) and, if appropriate, the acid acceptor are in general employed in approximately equimolar amounts. However, it is also possbile to use one component or the other in a relatively large excess (up to 5 mol).\nProcess (Excex1) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with carboxylic acid halides of the formula (VII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nSuitable diluents for the process (Eac) according to the invention are all solvents inert to the acid halides. Those preferably utilizable are hydrocarbons, such as benzine, benzene, toluene, xylene and tetraline, furthermore halogenated hydrocarbons, such as methylene chloride, chloroform, carbon tetrachloride, chlorobenzene and o-dichlorobenzene, and also ketones, such as acetone and methyl isopropyl ketone, furthermore ethers, such as diethyl ether, tetrahydrofuran and dioxan, moreover carboxylic acid esters, such as ethyl acetate, and also strongly polar solvents, such as dimethylformamide, dimethylsulphoxide and sulpholane. If the stability to hydrolysis of the acid halide permits, the reaction can also be carried out in the presence of water.\nSuitable acid-binding agents in the reaction of process (Es) according to the invention are all customary acid acceptors. Those preferably utilizable are tertiary amines, such as triethylamine, pyridine, diazabicyclooctane (DABCO), diazabicycloundecene (DBU), diazabicyclononene (DBN), Huinig base and N,N-dimethyl-aniline, furthermore alkaline earth metal oxides, such as magnesium oxide and calcium oxide, and also alkali metal and alkaline earth metal carbonates, such as sodium carbonate, potassium carbonate and calcium carbonate and also alkali metal hydroxides such as sodium hydroxide and potassium hydroxide.\nThe reaction temperature in the process (Excex1) according to the invention can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between xe2x88x9220xc2x0 C. and +150xc2x0 C., preferably between 0xc2x0 C. and 110xc2x0 C.\nWhen carrying out process (Excex1) according to the invention, the starting materials of the formulae (I-1-a) to (I-4-a) and the carboxylic acid halide of the formula (VII) are in general each used in approximately equivalent amounts. However, it is also possible to employ the carboxylic acid halide in a relatively large excess (up to 5 mol). Work-up is carried out according to customary methods.\nProcess (Excex2) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are each reacted with carboxylic anhydrides of the formula (VIII), if appropriate in the presence of a diluent and if appopriate in the presence of an acid-binding agent.\nPreferred diluents for the process (Excex2) according to the invention are those diluents which are also preferred when using acid halides. Otherwise, a carboxylic anhydride employed in excess may also simultaneously function as diluent.\nPossible acid-binding agents added in process (Excex2) are preferably those acid-binding agents that are also preferred when using acid-halides.\nThe reaction temperature in the process (Excex2) according to the invention can be varied within a relatively wide range. In general, the reaction is carried out at temperatures between xe2x88x9220xc2x0 C. and +150xc2x0 C., preferably between 0xc2x0 C. and 100xc2x0 C.\nWhen carrying out process (Excex2) according to the invention, the starting materials of the formulae (I-1-a) to (I-4-a) and the carboxylic anhydride of the formula (VIII) are in general each used in approximately equivalent amounts. However, it is also possible to employ the carboxylic anhydride in a relatively large excess (up to 5 mol). Work-up is carried out according to customary methods.\nIn general, a procedure is used in which diluent and excess carboxylic anhydride and the resulting carboxylic acid are removed by distillation or by washing with an organic solvent or with water.\nProcess (F) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with chloroformic acid esters or chloroformic acid thioesters of the formula (IX), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nSuitable acid-binding agents for process (F) according to the invention are all customary acid acceptors. Those preferably utilizable are tertiary amines, such as triethylamine, pyridine, DABCO, DBU, DBA, Hunig base and N,N-dimethyl-aniline, furthermore alkaline earth metal oxides, such as magnesium oxide and calcium oxide, additionally alkali metal and alkaline earth metal carbonates, such as sodium carbonate, potassium carbonate and calcium carbonate and also alkali metal hydroxides such as sodium hydroxide and potassium hydroxide.\nSuitable diluents for the process (F) according to the invention are all solvents which are inert to the chloroformic acid esters or chloroformic acid thioesters. Those preferably utilizable are hydrocarbons, such as benzine, benzene, toluene, xylene and tetraline, furthermore halogenated hydrocarbons, such as methylene chloride, chloroform, carbon tetrachloride, chlorobenzene and o-dichlorobenzene, additionally ketones, such as acetone and methyl isopropyl ketone, furthermore ethers, such as diethyl ether, tetrahydrofuran and dioxan, moreover carboxylic acid esters, such as ethyl acetate, furthermore nitriles, such as acetonitrile, and also strongly polar solvents, such as dimethylformamide, dimethyl sulphoxide and sulpholane.\nWhen carrying out process (F) according to the invention, the reaction temperature can be varied within a relatively wide range. The reaction temperature is generally between xe2x88x9220xc2x0 C. and +100xc2x0 C., preferably between 0xc2x0 C. and 50xc2x0 C.\nProcess (F) according to the invention is in general carried out under atmospheric pressure.\nWhen carrying out process (F) according to the invention, the starting materials of the formulae (I-1-a) to (I-4-a) and the appropriate chloroformic acid ester or chloroformic acid thioester of the formula (IX) are in general each used in approximately equivalent amounts. However, it is also possible to employ one component or the other in a relatively large excess (up to 2 mol). Work-up is carried out according to customary methods. In general, a procedure is used in which precipitated salts are removed and the reaction mixture which remains is concentrated by stripping off the diluent.\nProcess (G) according to the invention is characterized in that compounds of the formula (I-1-a) to (I-4-a) are in each case reacted with compounds of the formula (X) in the presence of a diluent and, if appropriate, in the presence of an acid-binding agent.\nIn preparation process (G), about 1 mol of chloromonothioformic acid ester or chlorodithioformic acid ester of the formula (X) is reacted per mol of a starting material of the formulae (I-1-a) to (I-4-a), at 0 to 120xc2x0 C., preferably at 20 to 60xc2x0 C.\nDiluents which may be added, if appropriate, are all inert polar organic solvents, such as ethers, amides, sulphones, sulphoxides, and also halogenoalkanes.\nDimethyl sulphoxide, tetrahydrofuran, dimethylformamide or methylene chloride are preferably employed.\nIf, in a preferred embodiment, the enolate salt of the compounds (I-1-a) to (I-4-a) is prepared by addition of strong deprotonating agents, for example sodium hydride or potassium tert-butoxide, the further addition of acid-binding agents can be dispensed with.\nIf acid-binding agents are used, customary inorganic or organic bases are suitable; sodium hydroxide, sodium carbonate, potassium carbonate, pyridine and triethylamine may be mentioned by way of example.\nThe reaction can be carried out at atmospheric pressure or at elevated pressure; it is preferably carried out at atmospheric pressure. Work-up takes place according to customary methods.\nProcess (H) according to the invention is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with sulphonyl chlorides of the formula (XII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nIn preparation process (H), about 1 mol of sulphonyl chloride of the formula (XII) is reacted per mol of starting material of the formulae (I-1-a) to (I-4-a) at xe2x88x9220 to 150xc2x0 C., preferably at 0 to 70xc2x0 C.\nProcess (H) is preferably carried out in the presence of a diluent.\nSuitable diluents are all inert polar organic solvents, such as ethers, amides, ketones, carboxylic acid esters, nitriles, sulphones, sulphoxides or halogenated hydrocarbons such as methylene chloride.\nDimethyl sulphoxide, tetrahydrofuran, dimethylformamide and methylene chloride are preferably employed.\nIf, in a preferred embodiment, the enolate salt of the compounds (I-1-a) to (I-4-a) is prepared by addition of strong deprotonating agents (for example sodium hydride or potassium tert-butoxide), the further addition of acid-binding agents can be dispensed with.\nIf acid-binding agents are employed, customary inorganic or organic bases are suitable; sodium hydroxide, sodium carbonate, potassium carbonate, pyridine and triethylamine may be mentioned by way of example.\nThe reaction can be carried out at atmospheric pressure or at elevated pressure; it is preferably carried out at atmospheric pressure. Work-up takes place according to customary methods.\nThe process (I) according to the invention is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with phosphorus compounds of the formula (XIII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nIn preparation process (I), 1 to 2, preferably 1 to 1.3, mol of the phosphorus compound of the formula (XIII) is reacted per 1 mol of the compounds (I-1-a) to (I-4-a) at temperatures between xe2x88x9240xc2x0 C. and 150xc2x0 C., preferably between xe2x88x9210xc2x0 C. and 110xc2x0 C., to give compounds of the formulae (I-1-e) to (I-4-e).\nThe process (I) is preferably carried out in the presence of a diluent.\nSuitable diluents are all inert polar organic solvents, such as ethers, carboxylic acid esters, halogenated hydrocarbons, ketones, amides, nitriles, sulphones, sulphoxides, etc.\nAcetonitrile, dimethyl sulphoxide, tetrahydrofuran, dimethylformamide or methylene chloride are preferably employed.\nAcid-binding agents which may be added, if appropriate, are customary inorganic or organic bases, such as hydroxides, carbonates or amines. By way of example, sodium hydroxide, sodium carbonate, potassium carbonate, pyridine and triethylamine may be mentioned.\nThe reaction may be carried out at atmospheric pressure or at elevated pressure; it is preferably carried out at atmospheric pressure. Work-up takes place according to conventional methods of organic chemistry. The end products are preferably purified by crystallization, chromatographic purification or by so-called xe2x80x9cincipient distillationxe2x80x9d, i.e. removal of the volatile constituents in vacuo.\nProcess (J) is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with metal hydroxides or metal alkoxides of the formula (XIV) or amines of the formula (XV), if appropriate in the presence of a diluent.\nPreferred diluents for process (J) according to the invention are ethers such as tetrahydrofuran, dioxan and diethyl ether or else alcohols such as methanol, ethanol and isopropanol, but also water. Process (J) according to the invention is generally carried out under atmospheric pressure. The reaction temperature is in general between xe2x88x9220xc2x0 C. and 100xc2x0 C., preferably between 0xc2x0 C. and 50xc2x0 C.\nProcess (K) according to the invention is characterized in that compounds of the formulae (I-1-a) to (I-4-a) are in each case reacted with (Kxcex1) compounds of the formula (XVI), if appropriate in the presence of a diluent and if appropriate in the presence of a catalyst, or (Kxcex2) with compounds of the formula (XVII), if appropriate in the presence of a diluent and if appropriate in the presence of an acid-binding agent.\nIn preparation process (Kxcex1), about 1 mol of isocyanate of the formula (XVI) is reacted per mole of starting material of the formulae (I-1-a) to (I-4-a) at 0 to 100xc2x0 C., preferably at 20 to 50xc2x0 C.\nProcess (Kxcex1) is preferably carried out in the presence of a diluent.\nSuitable diluents are all inert organic solvents, such as aromatic hydrocarbons, halogenated hydrocarbons, ethers, amides, nitriles, sulphones or sulphoxides.\nCatalysts may, if desired, be added to accelerate the reaction. The catalysts employed can very advantageously be organotin compounds, for example dibutyltin dilaurate.\nThe reaction is preferably carried out at atmospheric pressure.\nIn preparation process (Kxcex2), about 1 mol of carbamoyl chloride of the formula (XVII) is reacted at 0 to 150xc2x0 C., preferably at 20 to 70xc2x0 C., per mole of starting material of the formulae (I-1-a) to (I-4-a).\nPossible diluents optionally added are all inert polar organic solvents, such as ethers, carboxylic acid esters, nitriles, ketones, amides, sulphones, sulphoxides or halogenated hydrocarbons.\nDimethyl sulphoxide, tetrahydrofuran, dimethylformamide or methylene chloride are preferably employed.\nIf, in a preferred embodiment, the enolate salt of the compound (I-1-a) to (I-4-a) is prepared by addition of strong deprotonating agents (e.g. sodium hydride or potassium tertiary butoxide), the further addition of acid-binding agents can be dispensed with.\nIf acid-binding agents are employed, customary inorganic or organic bases are suitable; those which may be mentioned by way of example are sodium hydroxide, sodium carbonate, potassium carbonate, triethylamine or pyridine.\nThe reaction can be carried out at atmospheric pressure or at elevated pressure, preferably at atmospheric pressure. Work-up takes place according to customary methods.\nThe active compounds are suitable for controlling animal pests, preferably arthropods and nematodes, in particular insects and arachnida, which are encountered in agriculture, in forestry, in the protection of stored products and of materials, and in the hygiene field. They are active against normally sensitive and resistant species and against all or some stages of development. The abovementioned pests include:\nFrom the order of Isopoda, for example, Oniscus asellus, Armadillidium vulgare and Porcellio scaber. \nFrom the order of the Diplopoda, for example, Blaniulus guttulatus. \nFrom the order of the Chilopoda, for example, Geophilus carpophagus and Scutigera spec.\nFrom the order of the Symphyla, for example, Scutigerella immaculata. \nFrom the order of the Thysanura, for example, Lepisma saccharina. \nFrom the order of the Collembola, for example, Onychiurus armatus. \nFrom the order of the Orthoptera, for example, Blatta orientalis, Periplaneta americana, Leucophaea maderae, Blatella germanica, Acheta domesticus, Gryllotalpa spp., Locusta migratoria migratorioides, Melanoplus differentialis and Schistocerca gregaria. \nFrom the order of the Dermaptera, for example, Forficula auricularia. \nFrom the order of the Isoptera, for example, Reticulitermes spp.\nFrom the order of the Anoplura, for example, Phylloxera vastatrix, Pemphigus spp., Pediculus humanus corporis, Haematopinus spp. and Linognathus spp.\nFrom the order of the Mallophaga, for example, Trichodectes spp. and Damalinea spp.\nFrom the order of the Thysanoptera, for example, Hercinothrips femoralis and Thrips tabaci. \nFrom the order of the Heteroptera, for example, Eurygaster spp., Dysdercus intermedius, Piesma quadrata, Cimex lectularius, Rhodnius prolixus and Triatoma spp.\nFrom the order of the Homoptera, for example, Aleurodes brassicae, Bemisia tabaci, Trialeurodes vaporariorum, Aphis gossypii, Brevicoryne brassicae, Cryptomyzus ribis, Aphis fabae, Doralis pomi, Eriosoma lanigerum, Hyalopterus arundinis, Macrosiphum avenae, Myzus spp., Phorodon humuli, Rhopalosiphum padi, Empoasca spp., Euscelis bilobatus, Nephotettix cincticeps, Lecanium corni, Saissetia oleae, Laodelphax striatellus, Nilaparvata lugens, Aonidiella aurantii, Aspidiotus hederae, Pseudococcus spp. and Psylla spp.\nFrom the order of the Lepidoptera, for example, Pectinophora gossypiella, Bupalus piniarius, Cheimatobia brumata, Lithocolletis blancardella, Hyponomeuta padella, Plutella maculipennis, Malacosoma neustria, Euproctis chrysorrhoea, Lymantria spp, Bucculatrix thurberiella, Phyllocnistis citrella, Agrotis spp., Euxoa spp., Feltia spp., Earias insulana, Heliothis spp., Spodoptera exigua, Mamestra brassicae, Panolis flammea, Prodenia litura, Spodoptera spp., Trichoplusia ni, Carpocapsa pomonella, Pieris spp., Chilo spp., Pyrausta nubilalis, Ephestia kuehniella, Galleria mellonella, Tineola bisselliella, Tinea pellionella, Hofmannophila pseudospretella, Cacoecia podana, Capua reticulana, Choristoneura fumiferana, Clysia ambiguella, Homona magnanima and Tortrix viridana. \nFrom the order of the Coleoptera, for example, Anobium punctatum, Rhizopertha dominica, Acanthoscelides obtectus, Hylotrupes bajulus, Agelastica alni, Leptinotarsa decemlineata, Phaedon cochleariae, Diabrotica spp., Psylliodes chrysocephala, Epilachna varivestis, Atomaria spp., Oryzaephilus surinamensis, Anthonomus spp., Sitophilus spp., Otiorrhynchus sulcatus, Cosmopolites sordidus, Ceuthorrhynchus assimilis, Hypera postica, Dermestes spp., Trogoderma spp., Anthrenus spp., Attagenus spp., Lyctus spp., Meligethes aeneus, Ptinus spp., Niptus hololeucus, Gibbium psylloides, Tribolium spp., Tenebrio molitor, Agriotes spp., Conoderus spp., Melolontha melolontha, Amphimallon solstitialis and Costelytra zealandica. \nFrom the order of the Hymenoptera, for example, Diprion spp., Hoplocampa spp., Lasius spp., Monomorium pharaonis and Vespa spp.\nFrom the order of the Diptera, for example, Aedes spp., Anopheles spp., Culex spp., Drosophila melanogaster, Musca spp., Fannia spp., Calliphora erythrocephala, Lucilia spp., Chrysomyia spp., Cuterebra spp., Gastrophilus spp., Hyppobosca spp., Stomoxys spp., Oestrus spp., Hypoderma spp., Tabanus spp., Tannia spp., Bibio hortulanus, Oscinella frit, Phorbia spp., Pegomyia hyoscyami, Ceratitis capitata, Dacus oleae and Tipula paludosa. \nFrom the order of the Siphonaptera, for example, Xenopsylla cheopis and Ceratophyllus spp.\nFrom the order of the Arachnida, for example, Scorpio maurus and Latrodectus mactans. \nFrom the order of the Acarina, for example, Acarus siro, Argas spp., Ornithodoros spp., Dermanyssus gallinae, Eriophyes ribis, Phyllocoptruta oleivora, Boophilus spp., Rhipicephalus spp., Amblyomma spp., Hyalomma spp., Ixodes spp., Psoroptes spp., Chorioptes spp., Sarcoptes spp., Tarsonemus spp., Bryobia praetiosa, Panonychus spp. and Tetranychus spp.\nThe active compounds according to the invention are distinguished by a high insecticidal and acaricidal activity.\nThey can be used to particularly good effect for controlling insects which are injurious to plants, such as, for example, against the larvae of the mustard beetle (Phaedon cochleariae), against the larvae of the green rice leaf hopper (Nephotettix cincticeps) or against the caterpillars of the cabbage moth (Plutella maculipennis) (cf. the Use Examples).\nThe active compounds can be converted into the customary formulations, such as solutions, emulsions, wettable powders, suspensions, powders, dusting agents, pastes, soluble powders, granules, suspension-emulsion concentrates, natural and synthetic materials impregnated with active compound, and very fine capsules in polymeric substances.\nThese formulations are produced in a known manner, for example by mixing the active compounds with extenders, that is liquid solvents and/or solid carriers, optionally with the use of surface-active agents, that is emulsifying agents and/or dispersing agents and/or foam-forming agents.\nIn the case of the use of water as an extender, organic solvents can, for example, also be used as auxiliary solvents. As liquid solvents, there are suitable in the main: aromatics, such as xylene, toluene or alkylnaphthalenes, chlorinated aromatics and chlorinated aliphatic hydrocarbons, such as chlorobenzenes, chloroethylenes of methylene chloride, aliphatic hydrocarbons, such as cyclohexane or paraffins, for example petroleum fractions, mineral and vegetable oils, alcohols, such as butanol or glycol as well as their ethers and esters, ketones, such as acetone, methyl ethyl ketone, methyl isobutyl ketone or cyclohexanone, strongly polar solvents, such as dimethylformamide and dimethyl sulphoxide, as well as water.\nAs solid carriers there are suitable:\nfor example ammonium salts and ground natural minerals, such as kaolins, clays, talc, chalk, quartz, attapulgite, montmorillonite or diatomaceous earth, and ground synthetic minerals, such as highly disperse silica, alumina and silicates, as solid carriers for granules there are suitable: for example crushed and fractionated natural rocks such as calcite, marble, pumice, sepiolite and dolomite as well as synthetic granules of inorganic and organic meals, and granules of organic material such as sawdust, coconut shells, maize cobs and tobacco stalks; as emulsifying and/or foam-forming agents there are suitable: for example non-ionic and anionic emulsifiers, such as polyoxyethylene fatty acid esters, polyoxyethylene fatty alcohol ethers, for example alkylaryl polyglycol ethers, alkylsulphonates, alkylsulphates, arylsulphonates as well as albumen hydrolysis products, as dispersing agents there are suitable: for example lignin-sulphite waste liquors and methylcellulose.\nAdhesives such as carboxymethylcellulose and natural and synthetic polymers in the form of powders, granules or latices, such as gum arabic, polyvinyl alcohol and polyvinyl acetate, as well as natural phospholipids such as cephalins and lecithins, and synthetic phospholipids, can be used in the formulations. Further additives can be mineral and vegetable oils.\nIt is possible to use colorants such as inorganic pigments, for example iron oxide, titanium oxide and Prussian Blue, and organic dyestuffs, such as alizarin dyestuffs, azo dyestuffs and metal phthalocyanine dyestuffs, and trace nutrients such as salts of iron, manganese, boron, copper, cobalt, molybdenum and zinc.\nThe formulations in general contain between 0.1 and 95 per cent by weight of active compound, preferably between 0.5 and 90%.\nThe active compound according to the invention can be present in its commercially available formulations and in the use forms prepared from these formulations, as a mixture with other active compounds, such as insecticides, attractants, sterilizing agents, acaricides, nematicides, fungicides, growth-regulating substances or herbicides. The insecticides include, for example, phosphates, carbamates, carboxylates, chlorinated hydrocarbons, phenylureas and substances produced by microorganisms."}
{"text": "1. Field of the Invention\nThe present invention relates to a paper shredder with a pair of cutting rollers and individual strippers mounted on support rods. The strippers extend essentially from the lower and upper outer sides of the cutting mechanism between the cutting disks to the respective roller body. The strippers preferably have crowned stripper portions.\n2. Description of the Related Art\nThe configuration and arrangement of stripper elements for paper shredders operating according to the cutting roller principle have undergone intensive development in the past decade. In addition to the further development of the stripper elements arranged between the cutting disks, solutions have been proposed increasingly in recent years in which additional stripper elements are provided directly on the outer diameter of the cutting disks.\nThese solutions prevented that cut particles stuck between strippers and cutting disks, which over time resulted in increased friction and a higher ,contact pressure and frequently led to blockage of the cutting mechanism. However, these known solutions are more complicated and expensive and, when the installed power remained the same, a reduction of the available cutting power occurred."}
{"text": "In the electrical distribution field, a great number of electrical or electronic devices are mounted within a switchgear cabinet, the electrical or electronic devices including for example miniature circuit breakers, residual current circuit breakers, lamps, bell transformers, timer and push buttons etc. Each of such devices includes a mounting face which is so shaped as to match with an elongated bearing rail, for example, a so-called DIN rail which has a Ω-shaped cross-section. The DIN rail is fixed by screw or nuts to a supporting portion of a cabinet frame of a switchgear. The DIN rail has oppositely-disposed side walls provided, at both free edges thereof, with outwardly extending flanges. On the mounting face of each electrical or electronic device there is obtained a U-shaped recess having, on one boundary edge thereof, an inwardly turned lip and, on the opposite boundary edge, an inwardly protruding catch urged by a spring towards the inwardly turned lip. When a device has to be mounted on the DIN rail, the inwardly turned lip is put to engage behind one outwardly extending flange, and the spring-loaded catch makes a snap fit behind the other outwardly extending flange of the DIN rail. Subsequently, the electrical terminals of such a device have to be joined one by one directly to an elongated busbar extending alongside the DIN rail. These mounting and connecting operations are rather wearisome and time-consuming because have to be repeated several times, i.e. have to performed for each of the above mentioned electrical/electronic devices. Furthermore, if it occurs that all these devices have to be moved to a new location within the cabinet, each of them has first to be electrically disconnected and then mechanically decoupled from the DIN rail. Subsequently, the screws or nuts, which previously might be covered by the electrical devices, must be unscrewed thus enabling the DIN rail to be removed and be fixed to the new desired location within the cabinet. The electrical devices can thus be remounted on the DIN rail in the new location.\nIt would be desirable to overcome all the above mentioned drawbacks, by providing a clamping element which is simple and cheap to manufacture, and which enables a plurality of electrical/electronic devices to be quickly and easily connected within a switchgear cabinet in a versatile way."}
{"text": "An integrator is a device to perform the mathematical operation known as integration. Integrators are used in a number of settings to sample data over periods of time. Integrators may be employed by laboratory equipment, test equipment, medical devices, etc. A number of different integrator circuits are known in the art."}
{"text": "When a catheter is placed through the skin into an anatomic structure, a certain amount of physiologic concretion usually is found to collect on the outer surface of the catheter during the period that it dwells in the body. In blood vessels, this concretion is usually fibrin or clot. In the bladder, urinary substances (uric acid, urea, blood-components, biofilm) coat the outer surface. In chest tubes, blood and fibrin material adhere to the outer surface. In the intestinal tract, biofilm and succus entericus or stool frequently form a cast on the outer surface of the tube. In the trachea, the outer surface of the tube gets coated with mucous and respiratory substances that form a cast or biofilm on the outer surface.\nThus, when such a catheter is removed from the body, some of this material is often “peeled” off of the catheter surface, and if the catheter passes through the body wall, much of this detritus is left within the anatomic structure that has been cannulated. When this is blood clot in a vessel, it can lead to clot formation with a venous embolus from a vein (usually harmless, but can lead to thrombophlebitis) or an arterial embolus, which can be more serious if it obstructs an artery downstream The GI tract can usually expel the retained material without difficulty. Clot or fibrin masses in the pleural space can be reabsorbed or cause pleural adhesions which are of little clinical consequence. Urinary debris can cause problems if over 3-4 mm in diameter if it cannot be passed via the ureter or urethra, and such catheter debris can act as a nidus for bladder stone formation. With a minitracheotomy tube, the debris falls off into the patient's trachea, where it can stimulate a coughing episode that can increase the morbidity in emphysema or chronic obstructive pulmonary disease (COPD) patients, who are the usual candidates for use of such a tube.\nCOPD is a slowly progressive disease of the airways that is characterized by a gradual loss of lung function. In the U.S., the term COPD includes chronic bronchitis, chronic obstructive bronchitis, or emphysema, or combinations of these conditions. It represents the fourth leading cause of death in the U.S. 12.1 million adults ages 25 and older reported being diagnosed with COPD in 2001. About 24 million adults have evidence of impaired lung function indicating that COPD is underdiagnosed. About 119,000 adults ages 25 and older died from COPD in 2000. While the COPD death rate for females more than doubled between 1980 and 2000, and the number of deaths for females surpassed the number for males in 2000, the overall age-adjusted death rate for COPD remained higher for males in 2000. The age-adjusted COPD death rate was about 46 percent higher in males than females and 63 percent higher in whites than blacks. COPD is the fourth leading cause of death in the U.S. and is projected to be the third leading cause of death for both males and females by the year 2020.\nPhysicians who would normally use a COPD minitracheostomy ventilation technique may avoid using the minitracheostomy tubes because of the problem of mucous plugs falling into the trachea and leading to severe coughing episodes which leave COPD patients breathless and can actually be life threatening. The consequences of such solid material falling into the trachea is analogous to the feeling we all get when we aspirate a small food particle or sip of fluid. This set of problems of retained, solid components falling into a patient's trachea, blood vessels, bladder or urinary tract has led to a need to redesign these catheters. A need exists in the art for a design of a minitracheostomy tube in the lungs or a catheter in the blood vessel, bladder or urinary tract that can be easily inserted and then removed from the patients on a regular biweekly basis without causing severe problem of mucous plugs falling into the trachea and leading to severe coughing episodes, or debris from the catheter devices causing further complications in the patient."}
{"text": "1. Field of the Invention\nThis invention relates to integrated circuit (IC) devices. Embodiments of the invention are applicable to and suitable for electronic equipment such as television receivers, video tape recorders (VTR's) and audio tape recorders, whose internal circuits are digitally controlled.\n2. Description of the Prior Art\nVideo or audio equipment such as digital signal processing television receivers, VTRs and audio tape recorders have recently been commercialized. Most of such digital equipment uses internal bus lines common to a central processing unit (CPU), a memory, and other control circuitry. The CPU fetches control signal values for each circuit of the electronic equipment and stored in the memory, which may be a read only memory (ROM), during the normal operation and sends the fetched data to the appropriate circuit via the internal bus, which may comprise or include a so-called serial inter-integrated circuit (I.sup.2 C) bus, so as to control the operation of the internal circuit. In addition, the CPU may control the circuits in the equipment in direct response to signals supplied manually through an external keyboard or a remote controller.\nSuch an internal bus is conventionally either of the two-wire or three-wire type. A two-wire internal bus comprises a data transmission line and a clock transmission line. The two-wire type bus may utilize the communication system disclosed in European Patent Application No. 81201168.2 (0 051 332). On the other hand, a three-wire bus comprises a data transmission line, a clock transmisssion line, and a so-called ident signal line used to identify the nature of each data block coming over the data transmission line.\nIf standardized, centralized and simplified testing of the internal circuitry during assembly and servicing of such electronic equipment in which internal bus lines are installed were possible, an overall cost reduction including a reduction in manufacturing cost could be achieved."}
{"text": "1. Field of the Invention\nThis invention relates to humidity indicator cards and more particularly, to an irreversible humidity indicator card, which is covered on both sides, which utilizes a deliquescent material which does not contain a dye.\n2. Description of Prior Art\nWhen shipping or storing many types of materials, particularly electronic components, it is desirable to know whether those components have been exposed to a particular level of humidity which humidity could cause damage to those components. For example, electronic components can be damaged by exposure to low levels of humidity, even where the exposure is only for short periods of time. These same electronic components may also be contaminated and/or damaged by exposure to dust and other particulate matter.\nTo address the need for the detection of humidity levels within shipping or storage containers, humidity indicators have been developed. There are generally two types of humidity indicators. One of these humidity indicators reversibly changes color upon exposure to particular humidity levels. Such reversible humidity indicators typically utilize cobalt chloride as the humidity indicator material. It changes color when exposed to predetermined levels of humidity and returns to its original color when the humidity level drops below that predetermined level. These reversible humidity indicators are used to indicate the current condition of a desiccant and/or the current humidity level within the storage container.\nThe second type of humidity indicator is an irreversible humidity indicator. These indicators are designed to detect a predetermined level of humidity and provide a visual indication of whether components stored in the containers where these humidity indicators are used have been exposed to that predetermined level of humidity even for short periods of time and even if the level of humidity drops below that predetermined level when the components are checked at a later time. Large changes in humidity levels sometimes occur where storage containers are used in relatively warm climates where the moisture level in the air rises and falls dramatically depending upon the temperature of the surrounding air. Under these conditions, a reversible humidity indicator might fail to indicate the temporary presence of high humidity within a storage container even though such high humidity may be sufficient to cause damage to the components present in the storage container.\nOne of the first irreversible humidity indicator devices was disclosed in U.S. Pat. No. 2,214,354, which disclosed the use of a calcium chloride material which was mixed with a water soluble dye and deposited on a porous surface material, such as a sheet of absorbent paper. Upon exposure of the absorbent sheet to a predetermined humidity level, the calcium chloride material liquifies and releases the dye in liquid form. The dye is then carried by capillary action onto the porous surface of the absorbent paper, where it produces a permanent and irreversible dye mark. The aforementioned patent further describes various deliquescent agents which may be employed to show different humidity levels.\nTo maintain a consistently low humidity level, shipping containers and long term storage containers usually contain desiccant materials. These desiccant materials dehydrate the storage area and are intended to maintain the humidity level within that storage area at a predetermined level. These containers are periodically opened to recharge or replace the desiccant materials placed within the container and/or to check the level of humidity in the storage container. After replacing the desiccant material, the container is again sealed. In order to determine whether the humidity level in these storage containers has ever reached certain critical levels, irreversible humidity indicators are also often placed within the containers with the desiccant materials. These irreversible humidity indicators can be reviewed at the same time that the desiccant material is being checked to determine whether a harmful humidity level has ever been reached in the shipping container.\nA series of relative humidity indicators, each utilizing a different deliquescent salt, are disclosed in a series of patents which were issued during the 1940\"\"s and 1950\"\"s including U.S. Pat. Nos. 2,460,065, 2,460,066, 2,460,067, 2,460,068, 2,460,069, 2,460,070, 2,460,071, 2,460,072, 2,460,073, 2,460,074, 2,526,938, 2,580,737, and 2,627,505. In addition, some humidity indicator cards are capable of showing different levels of humidity on the same card by use of a series of different deliquescent agents that change color at varying humidity levels, as disclosed in U.S. Pat. No. 2,249,867.\nHumidity indicator sheets and cards which contain deliquescent salts and dyes have commonly been used to detect the relative humidity level present within storage containers. See for example, U.S. Pat. Nos. 2,249,867, 4,034,609, 4,150,570, and 4,854,160. Button-type humidity indicators or xe2x80x9cplugxe2x80x9d humidity indicators are also sometimes used with packaging material and are disclosed, for example, by U.S. Pat. Nos. 2,716,338, 3,084,658 and 4,050,307. Another device for monitoring humidity levels, particularly in poured cement, is disclosed in U.S. Pat. No. 3,680,364.\nA multiple layer, reversible humidity sensing device containing a reflective layer, which is useful in viewing the changes in color of a humidity indicator card is disclosed by U.S. Pat. No. 4,034,609.\nA reversible humidity indicator card contained within transparent, flexible sheet materials with an impermeable front layer is disclosed in U.S. Pat. No. 5,224,373. This humidity indicator card is specifically designed for utilization with electronic components. It is formed as a xe2x80x9cwindowxe2x80x9d in a barrier bag.\nA delayed action, irreversible humidity indicator card is disclosed in U.S. Pat. No. 4,793,180.\nAll irreversible humidity indicator cards known hitherto are based on combinations of deliquescent salts and water-soluble dyes. In order to prepare humidity indicator cards that react at various humidity levels, different combinations of deliquescent salts and dyes must be chosen. Only cards that show the same change in color at each chosen humidity level are acceptable to users of these cards. Otherwise, if there are varying color changes, it may be difficult to determine whether a particular humidity level has been reached. W. B. Abel: Chemical Maximum Humidity Indicator Update Report, BDX-613-1989 and U.S. Pat. No. 3,898,172 teach that only certain combinations of salts and dyes are useful for this purpose as the dye is quite soluble in the saturated salt solution that is formed upon deliquescence. In addition, the solubility and color of the dye must be independent of pH changes that may be attributed to the deliquescence of the salt. To ensure a proper shelf-life of the indicator card, the dye also must not react with the salt in any way (e.g. redox reaction, acid-base reaction). It is quite difficult to use either the same dye or different dyes with the same color and different salts over the entire humidity spectrum.\nMixing the individual salts with the dye is an additional required step for the production of the irreversible humidity indicators previously known. If the chosen salt and dye have different particle sizes, inhomogeneous distribution of the dye in the salt may occur and lead to inhomogeneous color and appearance on the indicating spot of the humidity indicator. This problem can be overcome by milling salt and dye together, but this is not possible for all saltxe2x80x94dye combinations, especially if the salt already holds water of crystallization.\nFurther, many of the deliquescent salts when they absorb moisture and melt can cause corrosion to the products stored in the shipping containers in which the humidity indicator cards are utilized.\nIn addition, the blotter paper used to form conventional humidity indicator cards sometimes sheds paper fibers and lint. Such fibers and lint may damage products that are sensitive to dust, such as electronic components.\nAccordingly, it is an object of this invention to produce an irreversible humidity indicator card which solves the problems present with conventional irreversible humidity indicator cards.\nIt is a still further object of the invention to disclose an irreversible humidity indicator card which is formed of a composite structure which includes a darkened blotter paper which is useful in showing a predetermined level of humidity.\nIt is a further object of the invention to disclose an irreversible humidity indicator card which does not use a dye with its deliquescent material.\nIt is a further object of the invention to disclose an irreversible humidity indicator card which does not produce paper fibers or lint when in use.\nThis invention is directed to an irreversible humidity indicator card comprising\nan intermediate carrier member, containing front and back sides and one or more holes passing through the member,\na water vapor permeable, first outer layer completely covering the front side of the intermediate carrier member,\na deliquescent material placed within the holes in the intermediate carrier member,\na colored absorbent sheet material, placed against the back side of the intermediate carrier member, which sheet material covers the holes in the intermediate carrier member, and\na second outer layer secured to the back side of the intermediate carrier member, which covers the colored absorbent material and covers the back side of the intermediate carrier member.\nPreferably, the colored absorbent material is a darkened blotting paper which absorbs the deliquescent material when the deliquescent material absorbs moisture and liquifies.\nIn another preferred embodiment, the water vapor permeable first outer material and/or the second outer layer further comprise materials with anti-static and/or electrostatic charge dissipative properties.\nThe invention is also directed to a process for the production of the above-referenced irreversible humidity indicator card comprising\npreparing an intermediate carrier member containing a front and back side and one or more holes passing through the member,\nsecuring a clear, water vapor permeable, outer layer to the front side of the intermediate carrier member,\nplacing a deliquescent material within the holes of the intermediate carrier member,\ncovering the holes of the carrier member on the back side of the intermediate carrier member with a colored, absorbent sheet material, preferably a blotter paper, and\ncovering the colored, absorbent sheet material and the back side of the intermediate carrier member with a second outer layer.\nAlternatively, the colored, absorbent sheet material may first be secured to the second outer layer prior to securing that second outer layer to the back side of the intermediate carrier member."}
{"text": "The invention relates to fiber optic telecommunication systems and more specifically to chromatic dispersion compensation in such systems.\nThe tendency of a pulse of light propagating through an optical fiber to broaden is a result of the fact that different wavelengths of light pass through the fiber at different speeds. This speed differential which causes the pulse to broaden is termed chromatic dispersion. Chromatic dispersion presents a problem in modem optical communication systems because the tendency of light pulses to broaden as they propagate down the fiber causes the closely spaced light pulses to overlap in time. This overlap can have an undesirable effect since it restricts how closely spaced the pulses can be. This in turn limits the data bandwidth of the optical fiber.\nThere are many characteristics of dispersion. First order dispersion is the rate of change of index of refraction with respect to wavelength in the fiber. First order dispersion is also referred to as group velocity. Second order dispersion is the rate of change of the first order dispersion with respect to wavelength. Second order dispersion produces the pulse broadening. Third order dispersion is the rate of change of broadening with respect to a change in wavelength. This is often referred to as the dispersion slope.\nSeveral solutions have been proposed to mitigate the effects of dispersion in transmission fibers. One technique involves the use of a compensating optical fiber having an appropriate length and which has a dispersion that is opposite to the dispersion characteristic of the transmission fiber. The result is dispersion in the transmission fiber is substantially matched and canceled by the total dispersion in the compensating fiber. While this technique offers a solution to the dispersion problem, it may be impractical in actual use because of the attenuation due to the required length of the compensating fiber. In such a case, the total transmission length of the fiber is significantly increased thereby increasing the signal attenuation in the fiber. Furthermore, it may be difficult to find a fiber of the desired length with the required dispersion properties.\nIt is also difficult to design a fiber having a changing index of refraction across the diameter of the fiber (the fiber index profile) that will compensate simultaneously for the second and third dispersion orders. It is even more difficult to control the material properties of such fibers even in the most accurate fabrication process necessary to produce such fibers. In addition, the process of fabricating the single compensating chromatic dispersion fiber is expensive and generally not practical.\nWhen a pulse of light is transmitted through an optical fiber, the energy follows a number of paths which cross the fiber axis at different angles. A group of paths which cross the axis at the same angle is known as a mode. Sometimes it is necessary to limit or control the number of modes used in a transmission system. The fundamental mode LP01 in which light passes substantially along the fiber axis is often used in high bandwidth transmission systems using optical fibers commonly referred to as single mode fibers.\nThe dispersion properties of high order modes have been investigated at length. There is a dependence of high order mode dispersion on wavelength and on the properties of the fiber. By properly designing the fiber index profile it is possible to make the dispersion slope be positive, negative or zero. It is also possible to make the magnitude of the dispersion be negative, zero or slightly positive. Using these two properties one can either control or compensate for the dispersion in any transmission fiber.\nSystems have been developed to take advantage of higher order modes to compensate for dispersion in a typical optical communication system. In such systems it has been necessary to first convert the lower order fundamental mode of the light to a higher order spatial mode. This is accomplished using longitudinal mode conversion.\nPrior art methods for mode conversion are known as longitudinal mode conversion and are based on introducing a periodic perturbation along the fiber axis. The length of each period and the number of periods in these longitudinal converters must be determined accurately according to the wavelength, the strength of the perturbation, and the modes involved. By constructing a longitudinal mode converter it is possible to achieve good efficiency in transferring the energy from one mode to the other in a limited spectral bandwidth. This spectral property has been used in Dense Wavelength Division Multiplexing (DWDM) applications in telecommunications for other applications. Unfortunately, this technique is accompanied by significant energy attenuation and it cannot be used over broad spectral bandwidths.\nAnother deficiency associated with longitudinal mode converters is related to the fact that after the conversion, only a single mode should be present in the fiber. It can be difficult to discriminate between desired modes and undesired modes having almost the same group velocities because unwanted modes can appear at the output of the converter. As the modes propagate, modal dispersion occurs and the pulse broadens. Generally, longitudinal mode converters introduce significant energy attenuation and noise. Therefore, a trade-off must be made between having broad-spectrum capability and the demand for converting the original mode to a pure, single, high-order mode.\nOne such longitudinal mode converter is discussed in U.S. Pat. No. 5,802,234. Here, a single mode transmission fiber carries the LP01 to a longitudinal mode converter. Before conversion in this system, however, it is necessary to couple the single mode transmission fiber to a multimode fiber while maintaining the signal in the basic LP01 mode. This coupling is typically difficult to achieve without signal degradation and any misalignment or manufacturing inaccuracies can result in the presence of higher order modes. It is desirable that only the LP01 mode propagate initially in the multimode fiber in order to avoid significant noise that degrades the system performance and typically such coupling results in the propagation of additional modes.\nThe present invention overcomes the disadvantages of longitudinal mode converters and previous attempts to control dispersion in a fiber optic system.\nThe present invention relates to an apparatus and method for chromatic dispersion compensation of optical systems. The apparatus and method make use of multiple chromatic dispersion compensation optical fibers. The number of orders of dispersion that can be corrected increases with the number of compensation fibers in the apparatus. Specifically, N orders of dispersion can be corrected by serially coupling N chromatic dispersion compensation optical fibers.\nIn one embodiment, the present invention features a chromatic dispersion compensation module which includes a first and second dispersion compensation fiber. The optical signal is dispersion compensated by each compensation fiber. In one embodiment, one compensation fiber compensates for first order chromatic dispersion and the other compensation fiber compensates for second order chromatic dispersion. In another embodiment, at least one of the compensation fibers is optical coupled to a transmission fiber.\nIn another aspect, the invention features a method of compensating for chromatic dispersion in an optical system which includes the steps of receiving an optical signal from a transmission fiber, injecting the optical signal into a first compensation fiber and dispersion compensating the optical signal. The method includes the additional steps of injecting the optical signal exiting the first compensation fiber into a second compensation fiber, additionally compensating the optical signal, and injecting the optical signal into a second transmission fiber."}
{"text": "Some vehicles are equipped with a closure panel, such as a lift gate, which is driven between an open position (position 2) and a closed position (position 1) using an electric drive system. Hold systems have been proposed to provide such vehicles with the capability of assisting the operator of the closure panel, in order to maintain a third position hold (or position 2) during opening and closing operations, so as to help counteract the weight of the closure panel itself. Without these hold systems, the closure panel may sag back down at the top end of the operational opening range due to the closure panel weight providing a closure torque greater than an opening torque provided by the electric drive system. Such proposed hold systems are, in some instances, complex and expensive and may not offer adequate failsafe modes (in the event of electric motor failure or loss of power) while at the same time maintaining adequate manual efforts by the operator. Also recognized is a need to provide a counterbalance mechanism that can offer efficient counterbalance force customization for different closure panel weights and configurations (e.g. differing centers of gravity), including the ability to accommodate for third position hold or stop and hold functionality of the closure panel.\nFurther disadvantages of current hold systems include bulky form factors which take up valuable vehicle cargo space, requirement to have additional lift support systems in tandem such as gas struts and other counterbalance mechanisms, unacceptable impact on manual open and close efforts requiring larger operator applied manual force at the panel handle, undesirable force spikes that do not provide for smoother manual force/torque curves, requirement to use vehicle battery power to maintain third position hold, and/or temperature effects resulting in variable manual efforts required by the operator due to fluctuations in ambient temperature.\nIt is recognized that constantly applied forces in a counterbalance mechanism can be problematic due to variations in the geometry and/or operator positioning during the complete raise and lowering cycle of a closure panel, including the ability to provide for third position hold where desired."}
{"text": "1. Technical Field\nThe invention relates to a device and a method for flanging, preferably roll-flanging, a component part. The device and the method preferably serve for hemming, particularly preferably roll-hemming. The invention can in particular be used in the construction of vehicles, preferably the construction of motor vehicles. Body parts, for example a door, a wheel arch or other part of an outer side wall, the front or rear of a motor vehicle, may in particular be considered as component parts to be reshaped.\n2. Description of the Related Art\nA roll-hemming device is known from WO 03/024641 A1, comprising three hemming, rollers which are used consecutively in a plurality of processing runs and, for this purpose, can optionally be moved into a processing position in which they roll off on a hemming web and bead the hemming web. In order to completely bead the hemming web around a hemming edge, the known device requires three processing runs, between which the next hemming roller in each case has to be moved into its processing position.\nIt is therefore an object of the invention to reduce the time and therefore cost of flanging, required for multiple-stage flanging, preferably roll-flanging."}
{"text": "1. Field of the Invention\nThe present invention relates to a composite cooking appliance consisting of a gas oven and a microwave oven, or a cooking appliance such as a gas oven, and it particularly relates to a temperature control device therefor utilizing gas combustion.\n2. Description of the Prior Art\nThe remarkable advance of semiconductor technology has resulted in sophistication of control circuits, miniaturization of such circuits by higher integration, and reduction of the costs of such circuits by mass-production, and these electronic control circuits have come to be also widely used in household electric equipment.\nThe technique using intelligence based on this electronic control is making rapid inroads into various heating apparatuses including electric ovens, microwave ovens, gas ovens, and combinations thereof.\nOne of the most important factors in cooking appliances based on gas combustion is exact temperature control. Since the quality of cooking depends on temperature, it is important that a preset temperature suitable for a given heating load be maintained accurately. Conventional cooking appliances using gas have temperature control means, most of which have been liquid expansion control methods and bimetals.\nA gas oven will be taken up by way of example. There is an arrangement wherein a main burner is fired to produce hot air, which is fed into the heating chamber of the oven to cook a heating load therein. The operation of a gas oven of this arrangement will be described with reference to FIGS. 4 and 5. When the user sets the knob at a certain temperature, two temperatures, upper and lower limit temperatures T.sub.1 and T.sub.2, are set by a control circuit. When cooking is started, the main burners start firing and eventually the upper limit temperature T.sub.1 will be reached at time t.sub.2. Then, the main burners stop firing, so that the temperature in the oven heating chamber gradually lowers until the lower limit temperature T.sub.2 is reached at time t.sub.3. Then, the main burners start firing again, repeating this cycle henceforth. If the preset temperature is changed to a higher one after passage of time t.sub.7 (the upper and lower limit temperatures being T.sub.3 and T.sub.4, respectively), the upper limit temperature T.sub.3 is reached at time t.sub.8, as shown in FIG. 4, and henceforth the same operation as described above is repeated.\nThus, the conventional gas oven is designed to control temperature by the on-off operation of the main burners to maintain the oven heating chamber temperature at a preset value, but this design has the following drawbacks.\nThe pressure of household gas differs with districts. Even in the same district, the gas pressure available for the gas oven installed in a home will always very owing to the turning on and off of gas in other homes or in another room in the same home. Thus, if the gas pressure drops below the normal value, this decreases the rate of heat generation by the main burners, thus requiring a longer time than usual to reach the preset temperature and hence a longer cooking time.\nFurther, since the ambient temperature of the air surrounding the gas stove differs greatly between midsummer and midwinter, the cooking time will be different in midsummer, particularly in a drafty room.\nFurthermore, because of the design to which turns off all of the main burners upon reaching of the upper limit temperature and turns on all of the main burners upon reaching of the lower limit temperature, the actual difference in temperature between the two extremes is as great as about 10.degree. C., showing that this design fails to attain the purpose of maintaining a temperature at a constant value. Furthermore, the need for frequently turning on and off the main burners entails the drawback of shortening the life of the control circuit system, particularly the relay."}
{"text": "The present application relates to software security and more specifically to systems and methods for adjusting and configuring security system behaviors, such as by adjusting enterprise software access and permissions.\nSystems for facilitating configuring and adjusting software security features are employed in various demanding applications, including user account login and authentication functionality for websites, systems for controlling user access permissions to enterprise databases, and so on. Such applications often demand configurable security systems that may be readily adjusted in accordance with the needs of a given computing environment.\nConventionally, software security mechanisms are often relatively static. For example, once a user has established an account or has otherwise been granted certain permissions to access and use computing resources, those permissions and account access capabilities seldom change unless a user changes job roles, closes an account, signs up for new services, and so on.\nHowever, the relatively static nature of the security systems and tedious methods for reconfiguring security system behaviors (e.g., methods involving opening new accounts, applying for new access privileges, and so on), may not account for potentially rapidly changing security needs of a given computing environment.\nAccordingly, in some instances, where less security or more security would be appropriate, a static level of security is provided. In certain cases, this may reduce user productivity, such as by requiring a user to repeatedly log in to software if the user has not interacted with the software for a predetermined time interval."}
{"text": "This invention relates to a window pane provided with a plurality of heating resistors and moisture-sensing probes. The invention is particular useful in operative association with a circuit arrangement for automatically controlling the heating circuit of the glass pane, i.e., in association with a circuit arrangement which automatically switches the heating circuit on when the glass pane becomes fogged and switches it off once the transparency of the pane has been restored. It is known that glass panes equipped with an electric heating system have many industrial and civil uses. In its use on vehicles (motorcars, airplanes, ships), the electric heating has the purpose of freeing the external surface of the glass pane from snow or ice and its internal surface from fogging caused by the condensation of water vapor upon it."}
{"text": "This invention relates to furniture that is easily assembled and disassembled, i.e., knockdown furniture, such as typewriter stands, desks, files, cabinets and the like. These types of furniture have parts which can be shipped disassembled or in a collapsed state so as to fit in a relatively flat container. Upon reaching their destination the furniture is put together.\nSome knockdown furniture is in the form of parts which are hinged to each other. This type of furniture is easy to assemble and the designs shown in U.S. Pat. No. 1,445,566 of Stoll and U.S. Pat. No. 4,280,744 of Nalcano are representative of this hinged furniture. The problem with hinged furniture is that it does not form a small package.\nAnother type of knockdown furniture is in the form of separate parts which must be connected to each other after shipment. U.S. Pat. No. 2,615,771 of Curtis, U.S. Pat. No. 3,838,902 of Tenani, U.S. Pat. No. 4,145,098 of Alexander and British Pat. No. 1,477,631 of Harvey et al. are representative. The problem with this type of furniture is that it requires fasteners and tools for assembly.\nIn the past, efforts have been made to provide furniture which can be purchased by the consumer in its disassembled state and then assembled without tools or fasteners. One example of this is disclosed in U.S. Pat. No. 4,317,416 of Baum et al. in which screw heads on one furniture part are received in raised slotted openings in mating furniture parts.\nAlthough the concept of knockdown furniture is old in the art, more often than not, it is difficult to assemble with or without tools and fasteners. This is because the parts are either too complex and intricate to assemble easily, or because deciphering the instructions may require an inordinate amount of time. Furthermore, fabricated furniture assembled without tools or fasteners is often unsatisfactory because it generally lacks the strength and rigidity required for every day use.\nIt would be advantageous to have knockdown furniture which is sufficiently stable for everyday use, can be assembled without tools and fasteners, and can be disassembled and shipped in a flat package."}
{"text": "1. Field of the Invention\nThis invention relates to a tube securing device for affixing medical tubing to the skin of a patient and, more particularly to a device for securing a thoracostomy tube to the skin of a patient.\n2. Description of Prior Art\nDuring the placement of larger medical tubes such as thoracostomy tubes, endotracheal tubes, and nasogastric tubes, it is important to secure the tube in position. It is particularly important that the tube remain stable and without movement in the longitudinal direction. Currently, these tubes are most commonly secured by wrapping tape around the tube and taping to the skin, or by suturing the tubing to the patient's skin. The difficulty occurs when wrapping tape or suture materials around a smooth tube. When a wet tubing is pulled on, the tape and suture material can slip or come undone allowing movement of the tubing and the resulting inability to secure the tube. Furthermore, it is cumbersome to remove tape or suture material from the tube when repositioning the tube without moving in a longitudinal direction.\nMany devices have been described for securing endotracheal tubes and nasogastric tubes in position with respect to some external reference point on the skin of a patient. Representative of such tube holders are U.S. Pat. Nos. 4,449,527; 5,060,645; 5,069,206; and European Patent EP356683 A. Such devices use complex locking mechanisms and often rely on tape or sutures to secure the tubing to the device. Moreover, these devices frequently include several parts and are, therefore, not unitary. They are often bulky and include a significant area of the device which lies on the face or mouth for stabilization. Such devices often provide reduction in horizontal movement at the expense of ease in placement. While such tube holders may be preferable in endotracheal tubes, which are utilized in a long term chronic setting, no simple device, unitary is currently available for securing a thoracostomy tube to the patient's body. The convention is to use a suture tie material attached to the skin incision around the tube to secure it. No device is currently available that can attach to the suture tie material and thereby secure the tubing.\nThe foregoing limitations in prior art tube holders with respect to securing or anchoring a thoracostomy tube in position so as to prevent vertical or longitudinal motion suggests a different approach to holder design is appropriate. What is truly required is a unitary device which can be made inexpensively and can be placed quickly around a thoracostomy tube in either an emergency or elective setting. Further, it would also be desirable if such a tube holding device provided means for facile removal from the tubing without injuring or moving the tubing."}
{"text": "1. Field of the Invention\nEmbodiments disclosed herein generally relate to an electronic device that is configured to perform a desired function based on input received from a user.\n2. Description of the Related Art\nThe popularity of portable electronics, such as smart phones, touch pads, PDAs, portable computers, wireless keyboards, wireless mice, wireless speakers, gaming controllers and portable music players, has increased dramatically in the past decade. As people have become more reliant on electronic devices, they have found more and more uses in the home, business and mobile environments.\nAs the development of software applications that run on today's electronic devices have improved, the usefulness and our reliance on these types of electronic devices have increased. To improve the usefulness of an electronic device, it is common to connect a first electronic device to one or more electronic devices, so that the first electronic device can perform, or better perform, functions that are typically outside the first electronic device's main function. For example, to improve a user's audio experience it is often desirable to link one or more portable speakers and an audio source, such as a music player or smart phone, together to provide a richer and enveloping audio experience. However, due to limitations in standard wireless communication protocols and device software, it is a non-trivial task to setup and control the communication between these electronic devices.\nTherefore, there is need for a method and apparatus that allows electronic devices to be automatically and seamlessly configured to allow the devices to rapidly provide desirable information to the user and to avoid the above-mentioned problems. There is also a need for a software application and a control method that allows an electronic device to be easily controlled by the delivery of simple input(s) received from a user."}
{"text": "The present invention relates to novel support and joiner means for shelving of the spaced rod type having vertically spaced rods forming an open and downwardly projecting lip at the front of the shelf.\nPrior art means for supporting and joining such shelves in end to end relation are commercially available but such means as exemplified in U.S. Pat. Nos. 3,598,064 and 3,765,634 are relatively expensive to manufacture because they comprise molded or die cast parts requiring the use of costly dies for their manufacture.\nAccording to the invention a simple means is utilized both for supporting and for joining together shelve units in end to end relation. The tooling required for fabricating such means is extremely simple and the cost of fabrication is minimal resulting in substantial economies in production without any reduction in function or strength."}
{"text": "1. Field of the Invention\nThis invention relates to vertically disposed backsplashes separately attached to horizontally disposed countertops such as in kitchens, and more particularly those made of solid plastics having a curved or coved interface at the approximate intersection of the backsplash and countertop. The coved interface being to provide an aesthetically pleasing and easy to clean interface between the two members.\n2. Description of the Prior Art\nCountertops with backsplashes which extend upward for at least a short distance above the countertop are desireable and common. Backsplashes most often are placed against a wall, but sometimes extend vertically upward from the countertop in locations which are not against a wall, to serve to terminate or designate the countertop work space. Countertops and backsplashes made of solid plastics, that is, non-laminated plastic solid surfacing material have become quite popular in recent years in kitchens, bathrooms, and offices of both residential and commercial buildings. The popularity of the use of plastic solid surfacing materials to define countertops and backsplashes is primarily due to the material having the qualities of high durability and ease of damage repair, ease of cleaning and sterilization, and beauty.\nSolid surfacing materials as they are known in the plastics industry, are non-foamed, non-laminated polymer based materials, or in short, solid plastics in sheet form useful for defining surfaces. Plastic solid surfacing materials are most often manufactured and sold in sheet form, typically in 30 inch wide.times.12 foot lengths in thicknesses between 1/4 through 3/4 inches. Such plastic solid surfacing materials are sometimes referred to by the general public as cultured marble. These polymer based solid surfacing materials are typically manufactured substantially of polyester or acrylic resins, or alloys thereof, depending on the manufacturer, and often contain smaller quantities of other components or additives both natural and synthetic to form desired color, visual pattern designs, and other desirable physical and visual characteristics in the plastics.\nPlastic solid surfacing materials are available from several U.S. manufacturers such as E.I. duPont de Nemours & Co., Inc. of Wilmington, Del. 19898 U.S.A., who market their polymer based solid surfacing materials under the trademark of \"Corian\". \"Corian\" is a trademark of duPont which refers to their solid surfacing polymer based material for use as a building material. \"Corian\" is a substantially rigid, non-foamed, non-laminated, non-coated solid material composed primarily of acrylic components. \"Corian\" is most often made and sold in sheet form. U.S. Pat. No. 3,847,865 issued Nov. 12, 1974 to R. B. Duggins and assigned to E. I. duPont de Nemours & Co., teaches one formula for making plastic solid surfacing material of the general nature of that referred to in this disclosure.\nAnother manufacturer of polymer based solid surfacing material is the Nevamar Corporation located at 8339 Telegraph Rd., Odenton, Md. 21113 U.S.A. The Nevamar Corporation markets their solid surfacing material under the trademark of \"Fountainhead\". \"Fountainhead\" is a substantially rigid, non-foamed, non-laminated, non-coated solid material composed of a polymer alloy comprised mostly of polyester components normally having therein a smaller percentage of acrylic components. \"Fountainhead\" is most often made and sold in sheet form.\nAnother manufacturer of polymer based solid surfacing material is the Formica Corporation, located at 155-T Rte. 46, W., CN-980, Wayne, N.J. 07470 U.S.A. The Formica Corporation sells their solid surfacing material under the trademark name of \"Surell\". \"Surell\", like \"Corian\" and \"Fountainhead\", is a dense rigid plastic most often made and sold in sheet form. \"Surell\" is a substantially rigid, non-foamed, non-laminated, non-coated solid material composed substantially of polyester components.\nDuPont, the Nevamar Corporation, and the Formica Corporation, and several other companies not specifically mentioned, whom produce very similar polymer based solid surfacing materials to one another, manufacture and sell polymer based solid surfacing building materials in sheet form intended for use as walling or countertops, and sometimes make and sell cast or heat-formed shapes made of the same polymer based materials useful as kitchen and bathroom lavatories.\nSome of the recognized advantages of using polymer based solid surfacing materials such as \"Corian\", \"Fountainhead\" or \"Surell\" over other available materials such as wood, metal, ceramic tile, and high pressure plastic laminates for countertops, exists in the fact that the material is a solid, polymeric non-laminated structure which the color or decorative color patterns extend completely therethrough, and therefore allow repeated refinishing of the material over many years using abrasives or sanding. If polymer based solid surfacing material does become stained, burned or scratched so deeply that the damage cannot be removed with a common household abrasive cleanser, the damage can be easily removed by light sanding with steel wool or fine sand paper, and this due to the fact that the material is solid, and the color or visual patterns extend completely therethrough. Furthermore, plastic solid surfacing materials typically have a high tensile strength, are quite hard, dense and rigid, and are resistant to chipping, cracking, splitting, warping, burning, and staining, all of which cannot be said about many other materials which could be used as substitutes therefore. Another attractive quality associated with polymer based solid surfacing materials such as those sold under the tradenames of \"Corian\", \"Surell\" or \"Fountainhead\", is the ease of adhesive bonding with available colored glues. Additionally, the polymer based sheets can be easily cut to size or otherwise shaped with mechanical material removal methods and tools using sawing and shaping tools such as router bits, power saws and shapers and the like, similar to those used to cut and shape wood.\nPolymer based solid surfacing materials such as \"Corian\", \"Surell\" or \"Fountainhead\" may be manufactured at a relatively low price to very closely resemble in texture and visually simulate marble, granite, and other natural stone products which have long been desired and used as building materials due to recognition of the durability and beauty of such natural substances.\nIt is the above stated qualities and advantages among others which make plastic solid surfacing material a good choice, and an increasingly popular choice for a countertop with backsplash.\nDue to excessive costs primarily related to the large number of different available colors from which consumers may now choose, and sizes which would be required to be made, shipped and stocked, countertops of plastic solid surfacing material are not widely manufactured with the backsplashes as integral one piece units. Although plastic solid surface materials have in the past been manufactured as countertops with the backsplashes as integral one piece units by either casting or heat forming, these countertops with integral backsplashes have for the most part been relatively short, narrow, straight lengths generally specifically for bathroom vanities. Due to the significant number of available colors and color patterns of solid surfacing material in sheet form, coupled with the frequent need to custom cut to size and shape countertops and backsplashes, the majority of solid plastic countertops with backsplashes are custom cut from sheet material and assembled in a countertop fabrication shop or at the installation site. This custom cutting and fabrication of countertops with attached backsplashes is particularly prevalent with kitchens and the like having large amounts of countertop space with the countertops and backsplashes intersecting one another at angles such as 90 degrees.\nWith countertops having attached backsplashes made of plastic solid surfacing material, one simple prior art system, and the most common one, of affixing the backsplash onto the countertop simply involves a butt-joint, which includes adhering the backsplash to the countertop at a right angle, thereby forming a right angled interior corner as may be ascertained from the attached prior art FIG. 1 drawing. The disadvantage of this butt-joint method being that the right angled corner is difficult to keep clean since dirt and water naturally accumulate along the seam, and the sharp angle (90 degrees) of the joint renders it difficult to get a sponge or like cleaning tool into the tight corner. Therefore, a visible joint is rendered more visible with the accumulation of dirt, distracts from the aesthetic appearance of the countertop, and may be unsanitary.\nOne solution to the difficult to clean right angled joint between the countertop and backsplash is to fashion a radius or coved interior corner joint between the two members which provides much more open space to access with a sponge for cleaning.\nA prior art method for creating a radiused or coved interior joint or intersection in plastic solid surface countertops with attached backsplashes, as is illustrated in the attached prior art drawing FIGS. 2A and 2B, is to adhesively affix a rectangular strip of plastic solid surfacing material along the right angled interior corner of the abutted backsplash and countertop, and then mechanically cut to remove a portion of the rectangular strip, such as with a router, leaving a radius or cove between the backsplash and countertop. This process involves utilizing a specialized jig or tool guide for holding the router at a forty-five degree angle to cut the radius. There is little room for error with this procedure, since routing the cove too deeply would cut into the backsplash or countertop, and too shallow a cove would require extensive sanding, therefore the installer must be quite skilled in this procedure. This procedure is quite time consuming, requires a high level of skill, and is therefore relatively expensive. Another disadvantage of this type of coved joint is that the longitudinal edges of the cove are quite thin, feathering out to a very thin edge so as to make a smooth transition or blend into the backsplash and countertop without leaving a highly visible seam. These thin feathered edges, particularly the feathered edge which lies in the horizontal plane of the countertop, can result in an area quite susceptible to wear. Even though the thin feathered edges at first may blend in quite well with the backsplash and countertop, they can eventually wear thin with repeated cleaning, particularly with abrasive cleansers. Chipping of the thin edges of the cove can eventually result. The chipped edges leave recesses which accumulate dirt and become harder to keep clean with time. Sometimes the adhesive shows through the chipped thin feathered areas, which is not very aesthetically pleasing when the adhesive is not an exact color match to the material, and the adhesive is seldom an \"exact\" color match. The showing of the adhesive through the thin feathered edges is particularly noticeable on the feathered edge which lies in the horizontal plane of the countertop because of the angle at which overhead lighting strikes the joint, and not so noticeable with the feathered edge lying in the vertical plane of the backsplash, again, because of the angle at which overhead lighting strikes the seam. As the thin feathered edges wear, more of the adhesive is exposed, and more dirt begins to collect and the joint becomes more visible with time. The problem with wearing of the thin feathered edges is most prevalent with the lower feathered edge which resides in the horizontal plane of the countertop, and not so much with the upper feathered edge residing in the vertical plane of the backsplash, as normally the horizontal plane receives far more vigorous cleaning.\nAnother prior art method of creating a coved interior joint or intersection on countertops with attached backsplashes of plastic solid surfacing material, as is illustrated in the attached prior art drawing FIGS. 3A and 3B, is to countersink the lengthwise bottom edge of the backsplash and an edge of an elongated rectangular strip into a rabbet (recess) formed along the upper lengthwise back edge of the countertop. A portion of the rectangular strip is then mechanically removed in the same manner as the first described prior art procedure to define a coved or curved joint. This process also involves utilizing a specialized jig or tool guide for holding the router at a forty-five degree angle. There is also little room for error in the depth of the routing of the cove, and therefore the installer must be quite skilled in this procedure. In this procedure, only the upper edge of the cove lying in the vertical plane of the backsplash is thinly feathered, and this has not been found to be a major problem. However, this method is very time consuming, requires a relatively high level of skill, and is also therefore quite expensive.\nIt is these shortcomings in countertops with separately attached backsplashes made of plastic solid surfacing material which the present invention overcomes. There is a significant need for an improved system of attaching a backsplash to a countertop made of solid plastics with the resultant joint between the two members being curved or coved for easy cleaning, and one which is very durable, inexpensive, and which the seam(s) is virtually invisible for many years."}
{"text": "Scanner systems for acquiring an image of an object, such as a printed circuit board, are well known in the arts of imaging and automated optical inspection. Some conventional scanner systems include sensors comprising a linear array of sensor elements. Other conventional scanner systems include sensors comprising a two dimensional array of sensor elements. Some systems employing two-dimensional array of sensor elements have been configured, for example, to operate in a time delay integration (TDI) mode of operation to acquire and image of an object. Other system employing a two-dimensional array of sensor elements have been configured to acquire a sequence of non overlapping images of an object.\nTDI systems are well known in the art of optical scanning. In such systems, a sensor array is scanned over an object, such as a printed circuit board, by moving either the array or the object in a direction perpendicular to the rows of the array. The scanning speed and an array clock are synchronized so that in each column of the array, multiple sensor elements in sequence capture light from the same point on the object. Charge is accumulated between rows as the sensor array passes over the object so that sensor signals in each column are summed for each point on the object, thereby providing an image of the object with enhanced signal/noise ratio.\nCommon TDI sensors are based on charge-coupled device (CCD) technology, which allows the sensor signals to be summed by transferring charge along each column of the sensory array such that newly accumulated charge is added to charge having accumulated in previous rows of the column. Other TDI systems based on photodiode arrays, such as CMOS sensor arrays, are also known in the art.\nU.S. Pat. No. 5,750,985 and U.S. Pat. No. 5,909,026, the disclosures of which are incorporated by reference, both describe sensors that can be employed in a TDI type arrangement.\nApplicants' copending U.S. patent application Ser. No. 10/141,988, filed on May 10, 2002 and entitled “Optical Inspection System Employing a Staring Array Scanner”, the disclosure of which is incorporated by reference, describes an inspection system employing a two dimensional sensor array."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor substrate and a method of manufacturing the same, and a semiconductor device utilizing such semiconductor substrate. More particularly, it relates to a semiconductor substrate for bipolar IC and BiCMOS and a manufacturing method thereof, and a semiconductor device utilizing such semiconductor substrate.\n2. Description of the Background Art\nFIG. 7 is a cross sectional view showing a buried diffusion epitaxial wafer formed by utilizing a conventional silicon substrate for bipolar IC and BiCMOS. Referring to FIG. 7, in the conventional buried diffusion epitaxial wafer, a buried diffusion layer 12 is formed on a top surface of a single crystal silicon substrate 13 at a predetermined interval. An epitaxial layer 11 of silicon is formed on the top surface of single crystal silicon substrate 13 and buried diffusion layer 12. A polycrystalline silicon layer 14 having a thickness of 1.0-2.5 .mu.m is formed on a rear surface of single crystal silicon substrate 13.\nBuried diffusion layer 12 is formed for decreasing collector resistance of bipolar IC or BiCMOS formed in epitaxial layer 11. Polycrystalline silicon layer 14 is formed for conducting gettering utilizing crystal defects included therein. Heavy metals such as Cu, Fe, and Au entered during manufacturing process are removed from electrically active regions of semiconductor elements by gettering.\nA method of manufacturing a conventional buried diffusion epitaxial wafer shown in FIG. 7 will be described.\nSingle crystal silicon substrate 13 is grown by Czochralski method. Czochralski method is disclosed, for example, in Semiconductor Silicon Crystal Technology (1989), Academic Press, Inc., by Fumio Shimura, pp. 129-131. FIG. 8 is an illustration of manufacturing of a single crystal by Czochralski method. Referring to FIG. 8, polycrystalline silicon (not shown) is put in a quartz crucible 51 and heated by a heater 50 to be turned to melted silicon 52.\nThen, the melted silicon 52 is brought into contact with a seed crystal 200 which is pulled up while being rotated. Accordingly, a single crystal silicon (silicon ingot) 100 having the same crystal axis as seed crystal 200 is grown. The pulling up speed of the conventional single crystal silicon 100 is about 1.0 mm/min. The interstitial oxygen concentration of single crystal silicon 100 has been controlled to 9-16.times.10.sup.17 (atoms/cm.sup.3) according to the old ASTM (old American Society for Testing and Materials).\nThen, single crystal silicon 100 formed as described above is sliced thin so as to form single crystal silicon substrate 13 as shown in FIG. 7.\nThereafter, polycrystalline silicon layer 14 is formed on the rear surface of single crystal silicon substrate 13 by the CVD method at the temperature of 650.degree. C. Buried diffusion layer 12 serving as a floating collector of the bipolar element is formed in a predetermined region on the top surface of single crystal silicon substrate 13 at a predetermined interval by, for example, the ion implantation method. Epitaxial layer 11 is grown on single crystal silicon substrate 13 and buried diffusion layer 12 at the temperature of about 1100.degree. C.\nConventionally, the buried diffusion epitaxial wafer is formed as described above. Bipolar IC or BiCMOS is formed in epitaxial layer 11 of such buried diffusion epitaxial wafer. For example, as shown in FIG. 7, a bipolar transistor including a base layer 16, an emitter layer 17, and a collector layer 18 has been formed on the main surface of epitaxial layer 11. An isolation layer 15 has been provided in order to isolate such bipolar transistors.\nConventionally, as described above, the interstitial oxygen concentration of single crystal silicon substrate 13 is controlled to 9-16.times.10.sup.17 (atoms/cm.sup.3) according to the old ASTM specification.\nHowever, in the conventional buried diffusion epitaxial wafer shown in FIG. 7, heat treatment is conducted at the temperature of 650.degree. C. at the time of formation of polycrystal silicon layer 14, and in addition, heat treatment at the temperature of 1100.degree. C. is conducted during the growth of epitaxial layer 11.\nIn this case, if interstitial oxygen concentration is as high as 15.0.times.10.sup.17 (atoms/cm.sup.3), generation of internal precipitation defects is accelerated because of heat hysteresis of the above-mentioned heat treatments. This results in crystal defects even in the element formation region of epitaxial layer 11 (a region up to a few .mu.m from the surface of epitaxial layer 11).\nWhen there is such a large amount of internal precipitation defects as to form crystal defects even in the element formation region, electrical characteristics of elements are deteriorated and warp of the buried diffusion epitaxial wafer is increased. Further, if crystal defects are precipitated even in the element formation region, crystal layers are slipped with each other. Such slipping is disclosed in, for example, Semiconductor Silicon Crystal Technology (1989), Academic Press, Inc., by Fumio Shimura, pp. 60-63, and pp. 286-289.\nIn the conventional buried diffusion epitaxial wafer shown in FIG. 7, gettering of heavy metal impurity is conducted by using crystal defects in polycrystalline silicon layer 14 and internal precipitation defects in single crystal silicon substrate 13. In this case, if single crystal silicon substrate 13 having a low interstitial oxygen concentration such as 10.times.10.sup.17 (atoms/cm.sup.3) is used, there is not provided enough internal precipitation defects to carry out gettering of the heavy metal impurity generated during formation of elements.\nWhen gettering is not sufficient, the heavy metal impurity is precipitated up to the surface of epitaxial layer 11 so that OSF (Oxidation Induced Stacking Fault) is formed on the surface of epitaxial layer 11 with the heavy metal impurity serving as a carnal. Such OSF deteriorates electrical characteristics of elements.\nConventionally, various problems as described above have occurred, because the range of the interstitial oxygen concentration of single crystal silicon substrate 13 is set rather wide without considering heat treatment and gettering effect during formation of polycrystalline silicon layer 14 and epitaxial layer 11.\nWhen forming single crystal silicon substrate 13 by Czochralski method, if the speed of cooling the single crystal silicon 100 which is pulled up is too fast, point defects generated during the growth of crystals are combined to form a large crystal defect.\nIf single crystal silicon substrate 13 shown in FIG. 7 is cut away from single crystal silicon 100 having such a large defect an inconvenience is generated. If there is a large defect in single crystal silicon substrate 13, many crystal defects are generated also on the surface of epitaxial layer 11 grown up on the top surface of single crystal silicon substrate 13. Thus, if elements are formed on the surface of epitaxial layer 11 including many crystal defects, the electrical characteristics of elements are deteriorated."}
{"text": "The present invention concerns a method for preparing unprotected hydroxy compounds or acylated derivatives thereof by conversion of silyl alkyl-protected hydroxy compounds. The invention also relates to novel intermediates useful in the method and for other purposes.\nWithin certain fields of synthetic chemistry, in particular within the field of preparation of glycosides, there exists a need for methods and intermediates making possible easy and high-yielding preparation of e.g. oligosaccharides and glycosides. Thus, the availability of e.g. suitably protected sugar \"building blocks\" and reactions capable of joining such \"building blocks\" on the one hand with one another and on the other hand with suitable aglycon units would greatly simplify the preparation of a wide variety of oligosaccharides, glycosides and glycoconjugates useful as i.a. synthetic biological receptors of the type described in Danish patent application No. 176/85 filed Jan. 14, 1985.\nRecently, Lipshutz et al, Tetrahedron Letters, 1981, 22, p. 4603, reported the preparation of 2-trimethylsilylethyl glycosides and their subsequent cleavage by LiBF.sub.4 to the corresponding free sugar using LiBF.sub.4 or a 1:1 mixture of LiF/BF.sub.3.Et.sub.2 O. However, Lipshutz does not indicate whether the cleavage product (or some intermediate) could have conserved its anomeric stereostructure."}
{"text": "It is well known that for transmission of a pictorial digital image over a narrow band communication channel, block transform coding techniques (such as discrete cosine transform DCT) achieve very high bit rate compression ratios. Spatial transform coding schemes are able to achieve such high bit rate compression ratios due to the fact that the spatial transformation concentrates most of the energy of a pictorial image in the transform coefficients representing lower spatial frequencies. When the transform coefficients are arranged in order of increasing spatial frequency, for example by the well known zig-zag scan technique, long runs of zero amplitude coefficients are generated. The values of the transform coefficients are Huffman encoded, and the strings of zero amplitude coefficients are run-length encoded for very efficients image compression. See U.S. Pat. No. 4,302,775, issued Nov. 24, 1981 to Widergren et al for an example of such a compression scheme in a video image compression system.\nUsing a block transform coding scheme, it is possible to transmit a full resolution television frame over a telephone line in less than one minute. It has also been suggested (see the article \"Image display techniques using the cosine transform\" by King N. Ngan, IEEE transactions on acoustics and signal processing, Vol. ASSP-32 No. 1, February 1984.) that in a low bit rate image transmission system it would be desirable to provide the viewer at the receiver with a quick (within several seconds) \"preview\" image during transmission of the the full resolution image, so that the viewer can quickly evaluate the usefulness of the image being transmitted, and terminate transmission if the full resolution image is not of interest, thereby effecting savings in time and cost. This capability would be very useful for example to quickly browse through a number of images (e.g..about.100) to find a particular image that is desired. Ngan suggests that such a low resolution preview image can be developed from the first few transform coefficients of an image by transmitting the transform coefficients in an order representing increasing spatial frequencies, and performing an inverse transformation on the first few coefficients that are received to provide a preview image.\nIn a practical pictorial image transmission system, the scheme suggested by Ngan for providing a preview image has several disadvantages which were not discussed in Ngan's article. In a block transform image compression scheme, the image is divided into blocks (e.g. 16.times.16) of pixels. Consecutive blocks of the image are transformed and the coded transform coefficients are transmitted. The progressive transmission scheme suggested by Ngan, wherein only the first few coefficients from each block are used to generate a preview image, was found by the present inventors to be impractical in an actual image transmission system, since transmitting only the first few coefficients of each block frustrates the run-length coding of the coefficients. Furthermore, since all the coefficients are needed to recover the full resolution image at the receiver, the lower frequency transform coefficients from all the blocks must be stored until the higher frequency coefficients are received, thereby increasing the memory requirements at the receiver.\nIt is the object of the present invention to provide a system for transmitting pictorial still images over a narrow band transmission channel such as a telephone line, including means for presenting a low resolution preview image that is free of the disadvantages noted above. It is a further object to provide such a system for transmitting a color pictorial image."}
{"text": "1. Technical Field\nThe present invention relates to an ink jet recording method and a recorded matter.\n2. Related Art\nThere is demand for an ink with a low viscosity in ink jet methods in which ink is discharged as minute droplets. On the other hand, because it is necessary to prevent bleeding of the ink on a recording medium, an ink jet ink that includes a penetrant is used (for example, refer to JP-A-2001-187850).\nHowever, in the recording method used in the related art, it is difficult to form an image at high speeds while sufficiently preventing bleeding of the ink depending on the type of medium."}
{"text": "The present invention relates to an incremental encoder and, more particularly, to an incremental encoder having a rotor provided with a plurality of indices and capable of detecting a reference or zero point without turning the rotor one full turn.\nEncoders have been widely used for electrically measuring distances and angles. Optical encoders and magnetic encoders are quite analogous in basic configuration.\nRotary encoders have been prevalently used for the electrical measurement of angles. In particular, optical encoders incorporate advanced optical techniques, they are capable of measurement with high accuracy and high resolution, they have high resolving power and are immune to disturbances such as external magnetism, and they have a long lifetime owing to their noncontact measuring operation.\nOptical encoders having such exemplary features are used in surveying instruments for measuring angles, for example.\nOptical encoders employed in current surveying instruments are classified into those of absolute type, i.e., absolute optical encoders, and those of incremental type, i.e., incremental optical encoders.\nIn an absolute optical encoder, angle and position on a circumference are in one-to-one correspondence, and positions on a circumference are registered as absolute addresses. Therefore, position information about any position can be obtained. However, an absolute optical encoder has a complex construction, and it is very difficult to construct an absolute optical encoder in a compact, lightweight manner for a surveying instrument.\nAs shown in FIG. 6, an incremental optical encoder comprises a rotor 9100 provided with a main scale 9110, a stator 9200 provided with a subscale 9210, and a detecting means having components 9310, 9320 and 9330. The main scale 9110 formed on the rotor 9100 has graduation lines arranged at equal angular pitch, e.g. of 80 sec, at the periphery of the rotor 9100 in an occulting pattern. The subscale 9210 formed on the stator 9200 has graduation lines arranged at an angular pitch equal to that of the main scale 9110 in an occulting pattern.\nThe detecting means comprises an LED (light emitting diode) 9310, a collimator lens 9320 and a photosensor 9330. The LED 9310 and the photosensor 9330 are disposed with the rotor 9100 and the stator 9200 therebetween. When the rotor 9100 rotates, the light emitted by the LED 9310 is interrupted once every time the rotor 9100 turns through an angle corresponding to one graduation of the main scale 9110, and the photosensor 9330 provides an electric signal indicating the interruption of the light. The electric signals provided by the photosensor 9330 are counted to detect the angle through which the rotor 9100 has turned.\nThe incremental optical encoder is able to start counting the output signals of the photosensor 9330 from any position of the rotor 9100 and to measure the angle of turn of the rotor 9100 from the position where counting the output signals of the photosensor 9330 is started.\nA surveying instrument employing an incremental optical encoder will be described with reference to FIGS. 7(a) and 7(b). A surveying instrument 10000 comprises a base unit 8100, a standard unit 8200 rotatably mounted on the base unit 8100 in a horizontal plane, a sighting telescope 8300 supported on the standard unit 8200 for turning in a vertical plane, a first detector 8400 for detecting the horizontal angle of the standard unit 8200, and a second detector 8500 for detecting the elevation angle of the sight line of the telescope 8300.\nThe base unit 8100 is connected by leveling screws 8160 to a leveling plate 8150 to be fixedly mounted on a tripod or the like. The level of the surveying instrument 10000 can be adjusted by turning the leveling screws 8160. The base unit 8100 is provided with a lower adjusting knob 8120 and a lower clamp knob 8130 to adjust and fix the base unit 8100. The standard unit 8200 is provided with an upper adjusting knob 8220 and an upper clamp knob 8230 to adjust and fix the standard unit 8200. The sighting telescope 8300 is provided with an elevation adjusting knob 8320 and an elevation clamp knob 8330 to adjust the elevation angle in the sighting direction of the sighting telescope 8300 and to fix the sighting telescope 8300 in an adjusted sighting direction.\nAn optical encoder included in the second detector 8500 detects the elevation angle of the sighting direction of the sighting telescope 8300 with respect to zenith, for example. The optical encoder of the second detector 8500 is provided with an index for zero detection to determine an angle from a reference. The signals are counted using the index as a reference or zero point to measure the angle. An optical encoder included in the first detector 8400 for measuring the horizontal angle of the standard unit 8200 has no particular reference direction and hence does not need any reference point. Therefore, the optical encoder is an ordinary one not provided with any index.\nThe incremental optical encoder provided with an index of the second detector 8500 will be described with reference to FIG. 8. The incremental optical encoder of the second detector 8500 comprises a rotor 8510, a stator 8520, and an optical detector comprising components 8531-8533 and 8535-8537 disposed with the rotor 8510 and the stator 8520 there between as shown. The rotor 8510 is provided at its periphery with a main scale 8511 having graduation lines formed at equal angular pitch, and a zero detection index 8512. The stator 8520 is provided with a first subscale 8521 for use in combination with the main scale 8511, and a second subscale 8522 for use in combination with the zero detecting index 8512. The optical detector 8530 comprises an index detecting unit and a main scale detecting unit. The index detecting unit comprises a first light emitting device 8531, a first collimator lens 8532 and a first photosensor 8533 and is capable of detecting the zero detecting index 8512 of the rotor 8510. The main scale detecting unit comprises a second light emitting device 8536, a second collimator lens 8535 and a second photosensor 8537. The main scale detecting unit detects light pulses produced by the occulting pattern of the main scale 8511 of the rotor 8510, and the second photosensor 8537 converts the light pulses into corresponding electric signals. The electric signals are counted to determine an angle from a zero detection point.\nA method of using the conventional surveying instrument 10000 thus constructed will be described. The leveling plate 8150 of the surveying instrument 10000 is mounted on a tripod, and the leveling plate 8150 is leveled by turning the leveling screws 8160. Then, a main switch is closed and the sighting telescope 8300 is turned one full turn to complete preparations. When the sighting telescope 8300 is turned, the rotor 8510 of the second detector 8500 is turned, and the zero detecting index 8512 sets a zero point to measure an angle from the zero point.\nThis conventional surveying instrument 10000 needs to turn by one full turn the sighting telescope 8300 provided with the second detector 8500 including the incremental optical encoder provided with the index; that is, the zero detecting index 8512 of the rotor 8510 must be detected by the index detecting unit of the optical detector 8530 by turning the sighting telescope 8300 one full turn. The sighting telescope 8300 needs to be turned by one full turn because the index 8512, in general, is invisible from outside. If the sighting telescope 8300 is turned quickly, the index detecting unit of the optical detector 8530 is unable to detect the zero detecting index 8512 of the rotor 8510 and hence measurement cannot be made. Accordingly, the sighting telescope 8300 must be turned at an appropriate turning speed below a certain turning speed, with the appropriate turning speed unavoidably being dependent on user intuition, which is very burdensome."}
{"text": "Valuable high purity lead compound and eventually metallic lead is obtained from minerals as well as from electrode paste or slime recovered from dismissed lead batteries by pyrometallurgy according either to a very high temperature (1500-1700° C.) process with attendant generation of large quantities toxic fumes to be abated, or to a high temperature (1000-1100° C.) process, as often practiced in case the starting material is recovered electrode paste, that implies a significant by-production of toxic slags to be disposed of as dangerous substances in special dumps.\nAn attempt to improve these situations has been pursued through a carbonation technique of the raw (impure) electrode paste in plants for producing secondary lead.\nThis technique gives substantial advantages in terms of reduced impact on the environment but poor economical results and therefore has not encountered diffused acceptance in the industry. Many have endeavoured to improve the performance of this technique as testified by the numerous patent publications such as: U.S. Pat. No. 5,827,347, also published as WO9966105; U.S. Pat. No. 5,840,262 also published as WO9858878; the European Patents No. 1619261, No. 1656463; the Italian Patent No. 01239001; the U.S. Pat. No. 4,336,236 and No. 1,738,081; the BP patents No. 239,257 and No. 272,053; the Belgian Patents No. 273,660 and No. 691,028 as well as the project “Cleanlead Gypsum” partially financiated by the European Union.\nLead is extracted from a lead mineral such as galena, through a high temperature (over 1500° C.) pyro-metallurgical process that produces impure lead. All sulphur present in the starting material is converted to SO2 such that generally to the pyro-metallurgical plant is associated a plant for producing sulphuric acid in order to avoid inertization of SO2 in the form of CaSO4, which has an extremely low commercial value."}
{"text": "1. Field of the Invention\nThe present invention relates to a vehicle seat in which a child seat is built into an occupant seat formed by a seat cushion and a seat back which is reclinable with respect to the seat cushion.\n2. Description of the Related Art\nThere are various types of conventional child seats. For example, there is a type in which the child seat is fixed to an ordinary occupant seat as a single unit. There is another type in which a child seat mounting/fixing member is built in an occupant seat so that the child seat is fixed to the occupant seat as a single unit. In yet another type, which will be referred to hereinafter as a \"built-in type\", a child seat is built in an occupant seat so that the vehicle seat can be changed between a state of use as an occupant seat and a state of use as a child seat as occasion demands.\nWith built-in type vehicle seats, there is no need to remove the child seat. Built-in type vehicle seats are therefore superior in that the only work involved therein is changing the seat from one state of use to the other (as an example, see Japanese Utility Model Application Laid-Open No. 2-23229).\nThe structure disclosed in the above-mentioned reference is as follows: in an ordinary occupant seat, a seat back and a connecting rod are swingable a predetermined amount with respect to a seat cushion. Vertical grooves are formed in a reverse surface of the seat back. A child restraining belt and a waist portion pad are accommodated in the vertical grooves. Accordingly, when the occupant seat is to be used, the vehicle seat is placed in its normal state. When the child seat is to be used, the seat back of the occupant seat is swung so that by using the seat cushion of the occupant seat as the seat cushion of the child seat as well, a child can sit facing the rear.\nHowever, in such a structure, the seat cushion of the occupant seat is used also as the child seat. Therefore, a drawback arises in that the angles of the seat surface are not structured so as to allow an optimal seat surface for a child. Further, in order to protect the sides of the child's body, it is necessary to form side wall portions. However, such side wall portions are unnecessary in an occupant seat. Namely, separate, independent cushion capabilities for an occupant seat and a child seat cannot be ensured."}
{"text": "1. Field of the Invention\nThe present invention relates to a garment pattern sizing template system and more particularly pertains to increasing and/or decreasing the size of garment patterns in a simplified manner.\n2. Description of the Prior Art\nThe use of pattern sizing tools of known designs and configurations is known in the prior art. More specifically, pattern sizing tools of known designs and configurations previously devised and utilized for the purpose of increasing or decreasing the size of garment patterns through known methods and apparatuses are known to consist basically of familiar, expected, and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which has been developed for the fulfillment of countless objectives and requirements.\nBy way of example, U.S. Pat. No. 1,214,296 to Gorton discloses an instrument for drafting patterns for garments. U.S. Pat. No. 4,672,748 to Perazzolo discloses a pair of tailoring set-squares for sketching the component parts of clothing especially for constructing garments of any size or model. U.S. Pat. No. 5,444,920 to Nelson discloses tools for use in dressmaking. Lastly, U.S. Pat. No. 5,570,533 to Vouyouka discloses an industrial pattern grading template.\nWhile these devices fulfill their respective, particular objectives and requirements., the aforementioned patents do not describe a garment pattern sizing template system that allows increasing and/or decreasing the size of garment patterns in a simplified manner.\nIn this respect, the garment pattern sizing template system according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in doing so provides an apparatus primarily developed for the purpose of increasing and/or decreasing the size of garment patterns in a simplified manner.\nTherefore, it can be appreciated that there exists a continuing need for a new and improved garment pattern sizing template system which can be used for increasing and/or decreasing the size of garment patterns in a simplified manner. In this regard, the present invention substantially fulfills this need.\nIn view of the foregoing disadvantages inherent in the known types of pattern sizing tools of known designs and configurations now present in the prior art, the present invention provides an improved garment pattern sizing template system. As such, the general purpose of the present invention, which will be described subsequently in greater detail, is to provide a new and improved garment pattern sizing template system which has all the advantages of the prior art and none of the disadvantages.\nTo attain this, the present invention essentially comprises a sheet of transparent plastic having limited flexibility. The sheet is in a circular configuration with a diameter of about 11 inches. The sheet has a center point with an aperture and a periphery. A plurality of radial construction lines extend from about the center point to about the periphery. A plurality of grids each have a center line with nine holes equally spaced along the center line for marking a pattern piece there beneath. Nine spaced lateral marking lines extend in alternating points from the holes with sequential marking numbers from a smaller number of size 6 to a larger number of size 14. A central hole is located on a radial construction line. A plurality of parenthetical primary positioning numbers on the sheet are located adjacent to the grids. A plurality of paper patterns each with a plurality of circled secondary positioning numbers correlate to the primary positioning numbers on the sheet for assisting a user in the proper positioning of the sheet on a pattern during operation and use. The radial construction lines include: a first line marked xe2x80x9c(2) back neck pointxe2x80x9d with an angled grid having the smaller number closer to the periphery and xe2x80x9c(1) center back neckxe2x80x9d with a transverse grid; a second line marked xe2x80x9c(4)a back waist and bottomxe2x80x9d with a radial grid having the smaller number closer to the periphery and xe2x80x9c(3)a back and side panel shoulderxe2x80x9d with an angled grid having the smaller number closer to the periphery; a third line marked xe2x80x9c(6) back sleeve balance notchxe2x80x9d with an angled grid having the smaller number closer to the periphery and xe2x80x9c(5) back shoulder endxe2x80x9d with an angled grid having the smaller number closer to the periphery; a fourth line marked xe2x80x9c(8) back waist and bottomxe2x80x9d with a radial grid having the smaller number closer to the periphery and xe2x80x9c(7) back top of side seamxe2x80x9d with an angled grid having the smaller number closer to the periphery; a fifth line marked xe2x80x9c(9) under arm seam topxe2x80x9d, xe2x80x9c(9)a under arm seam waistxe2x80x9d and xe2x80x9c(9)b underarm seam bottomxe2x80x9d with a radial grid having the smaller number closer to the periphery; a sixth line marked xe2x80x9cfront s/p (12) waistxe2x80x9d and xe2x80x9c(12)a bottomxe2x80x9d with a radial grid having the smaller number closer to the periphery and xe2x80x9c(11) front s/p shoulderxe2x80x9d with an angled grid having the smaller number closer to the periphery and xe2x80x9c(10) front s/p shoulder endxe2x80x9d with a transverse grid; a seventh line marked xe2x80x9c(14) front waistxe2x80x9d and xe2x80x9c(14)a front bottomxe2x80x9d with a radial grid having the smaller number closer to the periphery and xe2x80x9c(13) front mid shoulderxe2x80x9d with an angled grid having the smaller number closer to the periphery; an eighth line marked xe2x80x9c(16) lapel break @ neck xe2x80x9d with an angled grid having the smaller number closer to the periphery and xe2x80x9c(15) front neck pointxe2x80x9d with an angled grid having the smaller number closer to the periphery; a ninth line marked xe2x80x9c(18) front bottom of lapelxe2x80x9d with an angled grid having the smaller number closer to the periphery and xe2x80x9c(17) lapel @ collar notchxe2x80x9d with an angled grid having the smaller number adjacent to the periphery; a tenth line marked xe2x80x9c(19) front bottomxe2x80x9d with a radial grid having the smaller number closer to the periphery; an eleventh line marked xe2x80x9c(22) cuff t/slxe2x80x9d and xe2x80x9c(22)a cuff u/slxe2x80x9d with a radial grid with a larger number closer to the periphery and xe2x80x9c(21) sleeve head crownxe2x80x9d with an angled grid having the larger number closer to the periphery and xe2x80x9c(20) hind arm top and (20)a under sleevexe2x80x9d with an angled grid having the larger number closer to the periphery; a twelfth line marked xe2x80x9c(23) collar step end leftxe2x80x9d with an angled grid having the smaller number closer to the periphery and xe2x80x9c(24) collar step end rightxe2x80x9d with an angled grid having the larger number closer to the periphery; a thirteenth line marked xe2x80x9c(25) collar lapel break notch leftxe2x80x9d with a radial grid having the smaller number closer to the periphery and xe2x80x9c(25)a lapel break notch rightxe2x80x9d with a radial grid having the larger number adjacent to the periphery; and a fourteenth line marked xe2x80x9c(26) collar neck point notch leftxe2x80x9d with a radial grid having the larger number closer to the periphery and xe2x80x9c(26)a collar neck point notch rightxe2x80x9d with a radial grid having the larger number adjacent to the periphery.\nThere has thus been outlined, rather broadly, the more important features of the invention in order that the detailed description thereof that follows may be better understood and in order that the present contribution to the art may be better appreciated. There are, of course, additional features of the invention that will be described hereinafter and which will form the subject matter of the claims attached.\nIn this respect, before explaining at least one embodiment of the invention in detail, it is to be understood that the invention is not limited in its application to the details of construction and to the arrangements of the components set forth in the following description or illustrated in the drawings. The invention is capable of other embodiments and of being practiced and carried out in various ways. Also, it is to be understood that the phraseology and terminology employed herein are for the purpose of descriptions and should not be regarded as limiting.\nAs such, those skilled in the art will appreciate that the conception, upon which this disclosure is based, may readily be utilized as a basis for the designing of other structures, methods and systems for carrying out the several purposes of the present invention. It is important, therefore, that the claims be regarded as including such equivalent constructions insofar as they do not depart from the spirit and scope of the present invention.\nIt is therefore an object of the present invention to provide a new and improved garment pattern sizing template system which has all of the advantages of the prior art pattern sizing tools of known designs and configurations and none of the disadvantages.\nIt is another object of the present invention to provide a new and improved garment pattern sizing template system which may be easily and efficiently manufactured and marketed.\nIt is further object of the present invention to provide a new and improved garment pattern sizing template system which is of durable and reliable constructions.\nAn even further object of the present invention is to provide a new and improved garment pattern sizing template system which is susceptible of a low cost of manufacture with regard to both materials and labor, and which accordingly is then susceptible of low prices of sale to the consuming public, thereby making such garment pattern sizing template system economically available to the buying public.\nEven still another object of the present invention is to provide a garment pattern sizing template system for increasing and/or decreasing the size of garment patterns in a simplified manner.\nLastly, it is an object of the present invention to provide a new and improved garment pattern sizing template system comprising a sheet of transparent plastic. A plurality of radial construction lines extend from a region on the sheet at angles. A plurality of grids each have a central line with holes equally spaced along the central line and spaced lateral marking lines extend from the holes with sequential marking numbers from a lower number to a larger number with the central hole being located on a radial construction line. The radial construction lines from the second corner are sequentially marked with construction indicia.\nThese together with other objects of the invention, along with the various features of novelty which characterize the invention, are pointed out with particularity in the claims annexed to and forming a part of this disclosure. For a better understanding of the invention, its operating advantages and the specific objects attained by its uses, reference should be had to the accompanying drawings and descriptive matter in which there is illustrated preferred embodiments of the invention."}
{"text": "1. Field of the Invention\nThe field of invention relates to barbeque grill apparatus, and more particularly pertains to a new and improved barbeque grill housing apparatus wherein the same is arranged for mounting exteriorly of a balcony structure of an existing dwelling to afford protection to individuals relative to an associated barbeque grill, as well as permitting maximum utilization of space of the associated balcony structure.\n2. Description of the Prior Art\nBarbeque grill apparatus of various types is utilized throughout the prior art to permit outdoor cooking of various foods. A barbeque grill is frequently positioned within a balcony structure of an associated dwelling for convenience and access to the grill. Contemporary dwelling structure has limited available space of such balcony structure, wherein the instant invention attempts to provide a support housing to mount a grill exteriorly of a balcony structure, as well as affording limited access to children and the like preventing inadvertent injury relative to a barbecuing event. Barbeque apparatus is exemplified in the prior art in U.S. Pat. No. 4,836,480 to Besner illustrating a support frame for mounting a grill within the support frame.\nU.S. Pat. No. 3,915,146 to Bauer sets forth a grill apparatus including a housing, wherein the housing includes a gas filled container for directing the gas as fuel into the grill structure.\nU.S. Pat. No. 4,932,391 to Bierdeman sets forth a grill, wherein the grill includes a rotisserie structure for use with the grill structure.\nAs such, it may be appreciated that there continues to be a need for a new and improved barbeque grill housing apparatus as set forth by the instant invention which addresses both the problems of ease of use as well as effectiveness in construction and in this respect, the present invention substantially fulfills this need."}
{"text": "The invention relates to a depalletizing device for lifting up objects loaded onto a pallet. More in particular, the invention relates to a depalletizing device comprising a carriage as well as a friction roller attached to the carriage, which friction roller is rotatable about a horizontal rotation axis and is coupled to a rotational drive in order to rotate, the carriage being movable to and fro in a direction at right angles to the rotation axis of the friction roller in order to push the friction roller against a side of an object on the pallet in such a manner that the friction roller exerts an upward friction force on the object, as a result of which the object is raised on that side, and in order for the friction roller to pass underneath said object after the latter has been raised.\nGB 1 559 973 discloses a depalletizing device having a sliding carriage and a roller which is arranged at the front thereof so as to be rotatable. The roller is arranged in a spring-mounted manner relative to the carriage by means of a coiled spring. The roller is driven in rotation. During depalletizing, the carriage is moved in the direction of the objects on a pallet. When the roller touches an object on the pallet, the spring is pushed in, resulting in an increase in the spring force and thus the pressing force of the roller on the object. As a consequence of this increasing pressing force, the friction between the moving roller surface and the object increases until the roller moves the object upwards and causes it to tilt. One side of the object is lifted up by the roller and once the object has tilted sufficiently, the roller will pass underneath the object, the object being moved backwards via the roller."}
{"text": "1. Field of the Invention\nThis invention relates to ionizing radiation sensitive materials which are useful as electron beam resists and information recording media.\n2. Description of the Prior Art\nUnsaturated ester derivatives of a 1:1 copolymer of octadecyl vinyl ether-maleic anhydride exhibit electron sensitivity of about 10.sup.-7 to about 5.times.10.sup.-7 coulomb/cm.sup.2, but since the contrast of these materials is generally less than about 1.2 they do not have the ability to record submicron images. When exposed to electrons with energies below about 30 keV, other derivatives, namely the hydroxyethylpyrrolidone (HEP) esters of the same copolymer, have a sensitivity of only about 10.sup.-4 coulomb/cm.sup.2 to electron radiation.\nThe literature indicates that a mixture of electron sensitive polymers with widely differing sensitivities would be expected to have sensitivity intermediate between the values of the components of the mixture and to exhibit a greatly inferior contrast--and hence resolving power--to either of the constituents. I find that this is not the case with mixtures of the aforementioned hydroxyethyl pyrrolidone derivative with the mixtures of the products of reaction of the unsaturated ester derivatives. Even more significantly, it has been found that the contrast of such resists improves appreciably, which feature provides a resolution capability of less than 0.5 micrometer. This increased contrast is achieved herein with resists made by mixing the products of reaction of an alkylvinyl ether-maleic anhydride copolymer with a hydroxyalkyl acrylate such as hydroxyethyl acrylate, hydroxyethyl methacrylate, 3-hydroxypropyl acrylate, 3-hydroxypropylmethacrylate, 4-hydroxybutyl acrylate, 4-hydroxybutyl methacrylate or allyl or propargyl alcohol or pentaerythritol triacrylate and the products of reaction of an N-hydroxyalkyl amide or an N-aminoalkyl amide and an alkylvinyl ether-maleic copolymer anhydride to form the novel half-ester copolymer."}
{"text": "1. Field of the Invention\nThis invention relates to the design, construction, and use of a device for acute and chronic measurement of intracranial pressures in the human central nervous system, particularly during magnetic resonance (MR) imaging procedures, in particular during the injection or infusion of therapeutic drugs into the brain parenchyma. The present invention is an improvement over conventional methods and apparatus for measuring intracranial pressure, such as ventriculostomy with an external transducer, epidural strain gauge transducer, pressure sensitive epidural capsule, pressure sensitive subdural capsule, or subarachnoid pressure transducer.\n2. Background of the Prior Art\nElevated intracranial pressure (ICP) can result from a variety of pathophysiological disturbances and frequently presents a major problem in the management of patients with head injury or other neurological disorders and diseases, including brain tumors and intracerebral hemorrhage. Although numerous autoregulatory and compensatory mechanisms protect the brain against elevations in ICP, if the ICP is increased to levels approximately equal to or greater than those of the arterial blood pressure, cerebral blood flow can be severely compromised leading to brain herniation and death.\nThe primary aim of ICP monitoring is to provide an early warning of a deteriorating condition so as to enable corrective therapy to be implemented. The accuracy required when measuring ICP is typically of the order of 1-2 mm Hg, but equally important is the reliability of the data. To ensure reliability, the calibration of any measuring device must be easily checked and adjusted as necessary. In addition to the steady-state accuracy of any sensing device and recording system, it is essential to have adequate dynamic response characteristics in the measuring device.\nIt is also essential that measurements can be made with minimum interference to the patient. To satisfy this objective, non-invasive methods of determining changes in ICP indirectly have been attempted using, for example, impedance plethysmography. Such non-invasive methods have been used successfully in neonates. However, in adults, methods which give the most accurate measurements require access into the cranial cavity and are therefore more invasive.\nInvasive ICP devices have generally evolved in two basic directions. The first is based on implanting a sensor within the cranium. The second is based on mounting the sensor externally and connecting the measurement site through a fluid-filled transmission line. The three main sites for ICP Ring are the lateral ventricle, the extradural space, and the subdural or subarachmoid spaces.\nU.S. Pat. No. 4,014,319 to Favre discloses an intracranial pressure transducer comprising a small sealed capsule positioned in a trephined hole in the patient\"\"s skull, wherein a sensor diaphragm in contact with the dura mater is displaced by changes in intracranial pressure and produces an output signal proportional to the change in intracranial pressure. U.S. Pat. No. 4,026,276 to Chubbuck discloses a pressure monitoring apparatus implantable in the cranium, wherein the apparatus comprises a passive resonant circuit with inductance and capacitance capability for measuring intracranial pressure by comparison to a reference ambient pressure.\nU.S. Pat. No. 4,062,354 to Taylor et al. discloses an intracranial pressure transducer system comprising a holding bracket containing sensor elements which is positioned against the dura of the brain, wherein the elements within the holding bracket transmit electromagnetic signals related to the intracranial pressure to a receiver outside the patients body.\nU.S. Pat. No. 4,080,653 to Barnes et al. discloses a method and apparatus for recording intracranial pressure utilizing a transducer amplifier.\nU.S. Pat. No. 4,114,603 discloses an intracranial pressure monitoring device comprising a pressure-sensitive catheter insertable between the dura mater and arachnoid membrane. U.S. Pat. No. 4,114,606 discloses a monitoring apparatus for intracranial pressure measurement, wherein electromagnetic radiation is imposed on a passive circuit implanted in the cranium, the frequency at which the radiation is absorbed reflecting intracranial pressure.\nU.S. Pat. No. 4,147,161 to Ikebe et al. discloses a system for measuring or monitoring intracranial pressure which comprises a non-elastic detecting pouch inserted between the skull and the brain, wherein a pressure measuring device in the liquid in the pouch indirectly measures the intracranial pressure.\nU.S. Pat. No. 4,156,422 to Hildebrandt discloses an apparatus for treating hydrocephalus comprising housing which contains subcutaneously implantable components for measuring and controlling fluid pressure, wherein a second housing outside the patient contains measuring and control components whereby an intracerebral space may be automatically drained in response to a predetermined adjustable ICP.\nU.S. Pat. No. 4,210,029 to Porter discloses a differential sensor unit utilizing fiber optic light guides, wherein three light guides pass within a pneumatic line into one end of a rigid cylindrical envelope implanted in the skull. Detectors are arranged to actuate pressure display and pneumatic controls to adjust the internal pressure of the envelope to match the ICP and thereby measure the ICP. U.S. Pat. No. 4,265,252 to Chubbuck discloses an implantable transsensor device containing a passive RF resonant circuit which is responsive to changes in ICP.\nU.S. Pat. No. 4,385,636 to Cosman discloses an implantable telemetric differential pressure sensing device comprising a planar closed conductive loop which moves with a flexible diaphragm, wherein the resonant frequency of the conductive loop is detected telemetrically to determine pressure in a body compartment.\nU.S. Pat. No. 4,343,038 to McKean discloses a magnetic field generator and magnetic pick-up coil contained in an implanted ICP monitoring device. U.S. Pat. No. 4,354,506 to Sakaguchi et al discloses an intracranial pressure gauge which comprises a powerless resonance circuit composed of a coil and a condenser, a sensor equipped with an implantable pressure-sensitive section capable of changing the inductance or the capacitance of the condenser in response to a change in ICP. U.S. Pat. No. 4,438,773 to Letterio discloses an implantable subarachnoid bolt for use in measuring ICP. U.S. Pat. No. 4,465,075 to Swartz discloses an integrated circuit including a pressure transducer and temperature compensation circuit.\nU.S. Pat. No. 4,471,786 to Inagaki discloses a telemetering intracranial pressure transducer for detecting ICP, wherein a pressure-receiving layer disposed in contact with the dura mater transmits an output signal to an telemetering transmission circuit housed entirely beneath the patients scalp.\nU.S. Pat. No. 4,600,013 to Landy et al. discloses a probe for ICP measurements in the subarachnoid space, said probe comprising a threaded shaft having a lumen in contact with a pressure transducer.\nU.S. Pat. No. 4,621,647 to Loveland discloses an apparatus for monitoring and regulating ICP, wherein a manometer, transducer and regulator are interconnected by tubing and stopcocks. U.S. Pat. No. 4,627,443 to Chubbuck discloses an X-ray readable implantable pressure sensor,,wherein shifting of the radiopaque means is observed to indicate the change in pressure of a body cavity.\nU.S. Pat. No. 4,677,985 to Bro et al. discloses an intracranial probe to monitor both ICP and blood flow by thermal diffusion and hydrogen clearance techniques.\nU.S. Pat. No. 4,703,757 to Cohen discloses an optical fiber pressure transducer, wherein variations in transversely applied pressure to an elongated flexible fiber results in proportional light refraction and corresponding output signal.\nU.S. Pat. No. 4,711,246 to Alderson discloses a miniaturized pressure transducer, wherein light transmitted through a single optical fiber is reflected by a diaphragm in accordance with the pressure being measured.\nU.S. Pat. No. 4,723,556 to Sussman discloses an intracranial ventricular catheter assembly, wherein a pressure sensing device is connected to the proximal catheter. U.S. Pat. No. 4,738,267 to Lazorthes discloses an implantable intracranial pressure sensor, wherein a pressure-sensitive diaphragm placed directly on the dura mater transmits the intracranial pressure through a resistive type transducer. U.S. Pat. No. 4,772,257 to Hakim et al. discloses an external programmer for magnetically-adjustable cerebrospinal fluid shunt valve.\nU.S. Pat. No. 4,858,619 to Toth discloses an intracranial pressure monitoring system, wherein relief valves are combined with a pressure sensor for automatic venting of cerebrospinal fluid during elevated ICP.\nU.S. Pat. No. 4,805,634 to Ullrich discloses an adapter assembly for accurately positioning a removable biosensor implanted in the cranium.\nU.S. Pat. No. 4,841,986 to Marchbanks discloses a method and apparatus for measuring ICP, wherein a pressure sensor placed against the eardrum indirectly detects ICP changes based on displacement of the eardrum.\nU.S. Pat. No. 4,903,707 to Knute et al. discloses a ventricular catheter assembly comprising a catheter and a bolt, wherein the catheter can be inserted to a predetermined depth into the cranium. U.S. Pat. No. 4,995,401 to Bunegin et al. discloses a device for measuring anterior fontanelle pressure, wherein ICP changes can be detected noninvasively.\nU.S. Pat. No. 5,000,049 to Cooper et al. discloses an apparatus for measuring fluid pressures using a catheter device.\nU.S. Pat. No. 5,005,584 to Little discloses a fiber optic transducer utilizing a flexible membrane to transduce pressure by interrupting light transmission between fiber optic paths in a catheter or guide wire carrier.\nU.S. Pat. Nos. 5,074,310 and 5,117,835 to Mick disclose a method and apparatus for in non-invasively measuring changes in ICP , wherein a predetermined vibration signal is applied to one location in the skull and an output vibration signal is detected from another location in the skull reflecting ICP changes over time.\nU.S. Pat. No. 5,108,364 to Takezawa et al. discloses a monitoring catheter for medical use comprised of multiple tubes equipped for fluid delivery and removal, pressure measurement, and temperature measurement. U.S. Pat. No. 5,113,868 to Wise at al. discloses a pressure sensing catheter system comprising a catheter, a pressure sensor, and a signal conduit means within the catheter for signals between an external monitor and the pressure sensor, the signal conduit including two electrical conductors which are connectable to the external monitor.\nU.S. Pat. No. 5,191,898 to Millar discloses a cerebroventricular catheter means of measuring ICP and injecting or withdrawing cerebrospinal fluid, wherein a transducer positioned at the end of the cerebroventricular catheter and electrically coupled to a monitor allows for ICP monitoring.\nU.S. Pat. No. 5,292,899 to Watanabe et al. discloses a method for measuring ICP by using a device comprising a reservoir implantable under the skin of a patient and into which reservoir cerebrospinal fluid can be introduced from the ventricle of the patient, a flexible dome configured to be upwardly projected from said reservoir by the pressure of the cerebrospinal fluid and flexibly deformable according to an external force, a pressing part for pressing against said dome through the skin, a pressing-part-driving means for driving said pressing part pressing the dome and a flexible membrane provided at the tip of the pressing part and having an interior filled with a fluid, the method comprising: pressing the flexible membrane of the pressing part against the dome of the reservoir through the skin by means of the pressing-part-driving means after a zero point correction of the pressure transducer is performed by communicating the interior of the flexible membrane with the atmosphere.\nU.S. Pat. No. 5,325,865 to Beckman et al. discloses a catheter assembly for measuring fluid pressure in a body cavity, comprising an optical converter responsive to an electrical power source for energizing a light-emitting diode which has drift characteristics which vary in response to temperature, and a detection circuit. U.S. Pat. No. 5,230,338 to Allen et al. discloses an interactive image-guided system for displaying images corresponding to the placement of a surgical probe in the body. U.S. Pat. No. 4,173,228 to Van Steenwyk et al, and U.S. Pat. No. 5,042,486 to Pfeiler et al. disclose medical probes wherein electromagnetic signals are propagated between one antenna on the tip of the probe inserted into a body region and several antennae outside the body. The position and orientation of the probe tip are determined from the signals transmitted between said antennae.\nU.S. Pat. No. 5,211,165 to Dumoulin et al., U.S. Pat. No. 5,255,680 to Darrow and Dumoulin, U.S. Pat. No. 5,307,808 to Dumoulin et al., and U.S. Pat. No. 5,318,025 to Dumoulin et al. additionally disclose a tracking system in which radiofrequency signals emitted by an invasive device, such as a catheter, are detected and used to measure the device\"\"s position and orientation in a patient. Localization of devices in situ is achieved by transmit radiofrequency coils positioned at its distal end, which are detected by receive radiofrequency coils positioned around the imaging volume of interest. The position of the device, as determined by the tracking system, is superimposed upon independently acquired diagnostic images. U.S. Pat. No. 5,383,454 to Bucholz discloses a system for indicating a position of a tip of a probe which is positioned within an object on images of the object, wherein a computer employing translational software translate the position of the tip of said probe into a coordinate system corresponding to the coordinate system of the cross-sectional images.\nEach of the above cited patents provide advantages and disadvantages for ICP monitoring. For instance, intraventricular catheters are accurate but may be difficult to position in the presence of brain swelling and shifting of tissue. Subarachnoid devices are easily placed, but may underestimate higher ICP pressures, particularly if they are not coplanar to the brain surface. Epidural devices require placement through a burr hole and can be hampered by recording artifacts, dampening, and problems with calibration.\nMoreover, none of the above cited patents disclose a method or apparatus for evaluating changes in ICP resulting from the direct injection or infusion of drug agents into brain tissues. One of the significant problems with direct drug delivery to living tissue is assuring that the drug is accurately distributed to target receptor locations. The efficacious delivery of therapeutic agents for the treatment of brain tumors or neurodegenerative diseases, as two examples, requires that the agents be delivered as close to their receptors in the brain as possible, while minimizing increases in intracranial pressure during and after drug delivery. Liquid drug agents delivered into the brain through implanted catheters will disperse from the site of injection at variable rates depending on a number of factors, including the physicochemical characteristics of the drug, capillary uptake, metabolic degradation and excretion of the drug, size of the extracellular space (the volume fraction), and geometry of the brain cell microenvironment (tortuosity). The degree to which each of these factors influences the distribution of a particular drug agent may be an important determinant of the effectiveness of drug treatment of diseases of the central nervous system.\nIdeally, drug material injected into the brain infiltrates the extracellular space, and the subsequent tissue distribution of the drug is governed mainly by its molecular weight, molecular radius, and the tissue matrix structure into which the material has been injected. However, injection of a solution containing a macromolecular drug agent into the intraparenchymal extracellular space of the brain may instead result in the injected drug being sequestered as a cavity or depot. If the injected drug solution forms a fluid-filled cavity in the tissue, this may lead to tissue selling, an increase in ICP and, secondarily, altered interstitial transport of the drug solute.\nIt is therefore apparent that increases in ICP induced by intraparenchymal injections of liquid drug agents can injure tissues directly, or indirectly by retarding the efficacious distribution of the drug due to tissue swelling and retarded interstitial solute transport. Thus, it is important to be able to monitor any local and regional increases in ICP resulting from injections of liquid drug agents directly into the brain parenchyma. The availability of an MR-visible drug delivery device which incorporates a method means for monitoring local and regional changes in ICP would make it possible to obtain near real-time information on tissue pressure changes during interventional procedures in an intra-operative MR system.\nThe present invention discloses a device and method for monitoring intracranial pressure (ICP) during MR image-guided neurosurgical procedures, wherein an MR-compatible microsensor pressure transducer coupled to a pressure sensing diaphragm is positioned or located on the surface of an interventional device, such as at or near either a) a distal tip of the device and b) lateral surfaces of an MR-compatible device such as an MR-compatible multi-lumen drug delivery catheter. One or preferably more than one sensing element may be present on the interventional device. The microsensor and diaphragm detect and monitor local and regional changes in ICP during and after delivery (especially immediately after delivery while changes in local concentration of the injectant or drug agents can be observed) of liquid drug agents into the brain parenchyma. In at least one method within the scope of the invention, the pressure sensor includes a pressure transducer for converting a sensed pressure into a first electrical signal which is detectable by an external monitor transmission after transmission over a signal conduit means. There may be a first switching device for selectively controlling when the first electrical signal is generated, and a device that is responsive to at least one command signal from the external monitor for controlling the delivery of drug solution in response to detected alterations in ICP. In a practice of the present invention, the ICP monitoring device may also be inserted into any convenient location such as a lateral cerebral ventricle, cerebral cistern, subarachnoid space, subdural or extradural space, or venous sinus, under MR imaging guidance, and used to record intracranial pressures over seconds, minutes, hours to days in patients undergoing diagnostic or therapeutic neurologic interventions. Continuous monitoring of the ICP in real time may also be effected in this invention.\nThe invention also relates to an apparatus for the determination or monitoring of intracranial pressure comprising:\nan article that is insertable into the tissue or vasculature of the brain,\nsaid article comprising a pressure sensing element, a signal transmitting connector connecting the pressure sensing element and a signal reader (meaning artificial intelligence or artificial memory such as a reader, recorder or analyzer), the signal reader being able to (read, record, or analyze) signals from said pressure sensing element to indicate hydrostatic pressure or changes in hydrostatic pressure. The pressure sensing element may, for example, comprise a solid state pressure sensing element, a curved semiconductor element, or a physical sensor, alone or combined with electronic or electrical elements, such as a diaphragm (e.g., a diaphram or flexible, compressible, deflatable, inflatable or pressure moveable layer or element that has one surface that is exposed to an environment in which the pressure sensing element is placed, such as within the fluid or tissue fluid of a patient), as where a second surface of the diaphragm defines part of an internal cavity within the pressure sensing element, and movement of said diaphragm causes changes in pressure within the internal cavity. A pressure reading device may then be present within said internal cavity and the pressure sensing element is attached to the connector by electrical attachment of the pressure reading device to the connector. The pressure-sensing device may be connected to or independent of a device for the delivery of chemical agents into the bodily fluid, and particularly into the bodily fluids of the brain. Certain semiconductor elements are known to be pressure-sensitive, and therefore semiconductor elements may be positioned on the interventional device at the tip and/or at multiple locations on the lateral surface of the interventional device and be electrically connected to an outside source to interpret the signals.\nAt least one method according to the present invention for sensing the intracranial pressure within a patient comprises inserting the above described apparatus into the brain of a patient, visualizing the tissue around the apparatus and visualizing the apparatus itself by magnetic resonance imaging particularly in effect by continuous or real time visualization with MRI or fluoroscopy, less preferred), and positioning the apparatus within the device under magnetic resonance imaging, and reading signals from said pressure sensing element to indicate hydrostatic pressure or changes in hydrostatic pressure within the brain during or shortly after delivery of chemical agents into the fluids and/or tissue of the brain. Upon reading changes in hydrostatic pressure of the brain that exceeds a predetermined amount, an indication of the exceeding of that predetermined amount may be provided to an operator of the apparatus, as by an alarm, special image on a video screen, or by an automatic function change directed by memory (computer) associated with the system. For example, the apparatus may also comprise at least one lumen for delivery of chemicals to the patient, and chemicals are delivered to the patient through the at least one lumen, and upon reading changes in hydrostatic pressure of the brain that exceeds a predetermined amount during delivery of chemicals to the patient from the lumen, an indication of the exceeding of that predetermined amount is provided to an operator of the apparatus. This can be effected upon the indication of the exceeding of that predetermined amount by an artificial memory, that artificial memory then signaling the apparatus to reduce the amount of chemical being delivered to the patient.\nOne aspect of this invention is to provide a device and method for acute measurement of ICP in the human central nervous system during magnetic resonance (MR) imaging procedures, particularly those involving direct injection or perfusion of drugs into the brain parenchyma, with the measurement and sensing being performed during those stages of the drug delivery where pressure changes may occur.\nA second aspect of this invention is to provide a device and method for controlling the delivery of drug solution in response to detected alterations in ICP during magnetic resonance (MR) imaging procedures, particularly those involving direct injection of drugs into the brain parenchyma. This can be particulartly accomplished by the use of at least two pressure sensing elements, one distal and one proximal to the point of insertion of the interventional device. By having separated pressure sensing devices, a pressure gradient may be measured and/or the direction of flow and variations in flow from the point of introduction can be determined.\nA third aspect of this invention is to provide a device as an alternative method for the detection and measurement of intracranial pressure, replacing devices and methods such as conventional ventriculostomy with an external transducer, an epidural strain gauge transducer, a pressure sensitive epidural capsule, a pressure sensitive subdural capsule, or a subarachnoid pressure transducer.\nA fourth aspect of this invention is to provide a device and method to record intracranial pressures over extended time periods such as hours or days in patients that have been subjected to compression or shift of the cerebral ventricles resulting from severe head injury, tumor, or other causes of acute or chronic intracranial hypertension.\nA fifth aspect of the invention is to provide a device and method to accurately record ICP in the lateral ventricle, epidural, subdural, and subarachmoid spaces in patients in whom these fluid-filled cavities have been compressed by trauma or medical condition, including tumors or hematoma, or may have thick or mobile walls or small volumes that cannot be reliably penetrated with standard ventricular catheters when they are attempted to be positioned under stereotaxic guidance.\nA sixth aspect of this invention is to provide a device and method for safely measuring ICP with minimal risk of infection caused by meningitis or ventriculitis, minimal risk of intracranial hemorrhage from opening the dura or puncturing the ventricle, and minimal interference to the patient.\nA seventh aspect of this invention is to provide a device and method for MR image-guided therapeutic drainage of cerebrospinal fluid to relieve acute rises in ICP.\nAn eighth aspect of this invention is to provide a device and method for injection of a small volume of MR-visible gas for MR-ventriculography while monitoring ICP.\nA ninth aspect of this invention is to provide a device and method for a hydrocephalus pressure valve and self-sealing port for regulating the flow of cerebrospinal fluid through the catheter after placement of the catheter tip into the ventricle, subarachnoid, subdural, or epidural space under MR imaging guidance.\nA tenth aspect of the present invention is to provide a device substantially made of MR-visible elastomeric hydrogel, polymer, thermoplastic, or similar low friction material to minimize abrasive damage to the brain during insertion."}
{"text": "Generally, an iron and steel plant produces slabs according to the customer requirements on steel products, such as specification, steel grade, chemical composition, hardness and so on. During the steel production process, first of all, a production schedule is made for steelmaking and continuous casting according to the production order. After slabs are produced, they are subjected to further treatments such as hot-rolling, cold-rolling and so on according to the customer requirements on the products. Finally, finished products are produced to be delivered to the customers. However, the following several reasons cause a low slab utilization, leading to a lot of irrational slab-order matching relationships: 1) In the current iron and steel market, orders tend to require a large variety and a small volume production, which is contradictory to batch production mode of iron and steel industry, for example, a converter or an electric furnace can handle 300 or 250 tons of molten steel in a batch, such that the amount of slabs in accordance with the order may be more than that required by the order; 2) when allocating the open-order slabs to orders, the matching relationship between the open-order slabs and the orders may be sub-optimal due to lack of working experiences and complete information; 3) the generation of unqualified products in post operations may cause that the amount of the products remaining for the order, to which the unqualified products belong, is less than that required by the order. Then, some more slabs need to be allocated to this order.\nIn prior literatures, most of the domestic and foreign studies on the slab matching problem aim at open-order slabs, which is about the matching of the surplus inventory, while the present invention tries to re-optimize the matching relationship in consideration of the original slab-order matching relationship and the load balance of all equipment units, so as to improve the slab-order matching relationship, decrease the slab cutting loss, reduce the redundant slabs weight of each order, and improve the integrity of an order, thus improving slab utilization, ensuring to deliver an urgent order on time, and improving the customer satisfaction."}
{"text": "U.S. Pat. No. 4,155,441 to Albrecht et al. discloses a conveyor system for separating columns of articles into spaced groups, each group containing a predetermined number of articles. The columns of articles are fed onto a conveyor having an extendable noser arrangement at the output end that is positioned over the input end of a second conveyor. The noser arrangement is extended to prevent the flow of articles from the first conveyor to the second, faster conveyor, thereby providing a space on the second conveyor between successive groups or articles. During extension of the noser arrangement, no articles are transferred since the noser arrangement extends the conveyor at the same rate that the articles are carried toward the end of the conveyor. A photoelectric device counts the articles as they are transferred and actuates the noser extending mechanism. U.S. Pat. No. 4,155,441 does not disclose changes in the noser arrangement due to the detection of an unevenly spaced article. In addition, this patent does not disclose continuous correcting of unevenly spaced articles during transfer to a second conveyor belt.\nDevices for shingling articles are known. For example, FIG. 1A shows a side view of a typical shingle stacking conveyor arrangement. The speed of the first belt (V1) is faster than the speed of the second belt (V2). This causes the articles on the second belt to be spaced closer together and to overlap or shingle. FIG. 1A shows an article that is ahead of where it would be if it were evenly spaced, causing it to arrive early at the second conveyor. FIG. 1B shows an article that is behind the position it would occupy if it were evenly spaced, causing it to arrive late at the second conveyor.\nThe effect of uneven article spacing on the first conveyor is to cause an uneven article spacing on the second conveyor, which results in an uneven shingle pattern (i.e., the articles are shingled at different angles). If the product arrives sufficiently ahead or behind of an evenly-spaced position, then the shingle pattern can be lost altogether, resulting in articles being piled on top of each other in a random manner, as shown in FIG. 1C.\nRandomly spaced or randomly stacked products is not just a matter of aesthetics. Unevenly distributed product does not go through conveyor systems to wrapping systems well. Any disruption in the pattern of products only gets worse as it makes its way to wrapping. Loss of product overlap, buckled product, and product jams all occur with regularity. When the disordered product arrives at the wrappers, it is prone to jam in chutes, feed mechanisms, infeed track, and the like. As a result, serious product breakage and wrapper downtime occurs. The present invention provides an apparatus that can take a single row of unevenly spaced articles on a first conveyor and ensures that the articles are more evenly spaced or more evenly stacked on a second conveyor."}
{"text": "The present invention relates to electrostatic filters in general, and more particularly to improvements in so-called high tension (ionic bombardment) filters. Still more particularly, the invention relates to improvements in a method and apparatus for automatically regulating the operation of high tension filters by regulating the potential which is applied to such filters.\nIn presently known high tension filters, the potential which is applied thereto is increased to reach the breakdown value and is thereupon reduced to a variable extent and for a variable interval of time to a value below the preceding breakdown value. Such operation is followed by a renewed increase of potential to the breakdown value. This is deemed to be advisable and advantageous because the electrical filter output can be regulated in a more satisfactory way to follow the varying breakdown resistance of the gaseous carrier medium for solid impurities which require segregation from the carrier medium.\nThe operating potential of a high tension filter is invariably limited by spark discharge between the corona discharge electrode and the collecting electrode of the filter. As a rule, the filter potential is selected in such a way that some arcing in the filter will take place because the rate of separation (i.e., the separation efficiency) is then at a maximum value. On the other hand, the frequency of arcing should not be too high.\nThe aforediscussed prior filtering methods and apparatus exhibit the drawback that each arcing leads to a total collapse of the electric field and, by using modern operational switches (thyristors), each arcing is followed by a complete shutdown of the supply of potential for a period of a few half waves in order to avoid the initiation of an immediately ensuing follow-up arcing. This entails the development of breakdown times during which the charging does not take place in an optimum way and to an interruption of field forces which are required for separation."}
{"text": "1. Field of the Invention\nThe present invention relates to a fluid control system, and in particular to a fluid flow control valve. In one application, the fluid flow control valve can check the free flow of water from a broken or damaged riser of an overhead irrigation system.\n2. Description of Related Art\nOverhead irrigation systems often incorporate water spray devices (e.g., sprinkler heads) mounted on risers. A riser supports a sprinkler some distance above the ground, e.g., one to three feet high. At this elevated position, the sprinkler has an enlarged spray pattern and irrigates a larger area than if the sprinkler were positioned closer to the ground. Riser mounted sprinklers also accommodate crops of varying heights.\nRisers are pipes or conduits, typically made of plastic such as polyvinylchloride (\"PVC\"), copper or galvanized steel. In many irrigation systems, the riser extends upward from a fitting of the irrigation system, such as from a \"T\" or an elbow juncture, located under ground. In such an arrangement, water flows from a subterranean irrigation pipe through the riser to the sprinkler.\nWhile the use of a riser-mounted sprinkler head enlarges the area irrigated by the sprinkler, the riser is susceptible to mechanical damage because it extends above ground in an exposed position. No external structure typically braces or buttresses the riser. Plastic or metal risers are thus easily broken (e.g., severed) or otherwise damaged (e.g., theft, vandalism, etc.), resulting most often in an open, free-flowing outlet.\nWhen a riser is broken or the sprinkler is removed, water cannot be properly distributed through the sprinkler head. The rate of water flow increases without the restricting back pressure provided by the sprinkler, such that a large stream of water projects above the broken riser. The resulting water geyser impacts against a relatively small ground surface. Serious flooding and erosion consequently results in a small area, while the remainder of the area normally irrigated by the sprinkler goes unwatered. A significant amount of water is wasted as the result of the unrestricted flow through the broken riser, and substantial soil erosion can occur.\nIn addition, the water fountain gushing from the broken riser also can pose a serious highway problem if the water sprays onto highway lanes, or if the resulting water and soil run-off flows onto the highway. Numerous automobile accidents occur each year due to broken irrigation system risers.\nWhile valves exist in the prior art for control of fluid flow in general, many such valves are inappropriate for restricting the flow of irrigation water through broken risers. Additionally, high volume farm irrigation systems introduce special needs for flow control. Prior valves are overly complicated and expensive for application to farm irrigation systems, which utilize hundreds of riser-mounted sprinklers. Some devices which have been implemented for stanching the unrestricted flow of irrigation water have proven unreliable, often failing to stop or even slow the rate of water flow when a riser breaks or other failure occurs. Other prior devices are too sensitive, shutting down water flow to undamaged risers.\nIn addition, when prior devices do function to stop the flow of water through a broken riser, the broken sprinkler often remains undetected for days, leaving the area surrounding the broken sprinkler unwatered. This danger is especially true of farm irrigation, where a single broken riser in a large field with many sprinklers could easily escape notice for many days, damaging crops in the vicinity.\nFurther deficiencies in prior devices include an inability of these devices to cope with transient flow conditions. For instance, the prior devices may prematurely shut off the valve in response to the combination of air and water in the system that often occurs when the system is first turned on. In addition, prior devices often include improperly restrained components which vibrate and wear under normal flow conditions.\nPrior devices also have close fitting parts which are subject to corrosion. As a result, the small space between the parts cannot be maintained as corrosion, scale, debris, etc., often fill the space between the closely fit parts, and the device is likely to malfunction. For example, in a prior valve, a disc or poppet is shaped and sized to have a close fit with the valve seat in order to arrest fluid flow under abnormal flow conditions. Upon corrosion, however, the disc or poppet will not properly seal the valve seat, thereby allowing fluid flow within the valve. Under some conditions, such corrosion, scale formation, or like collected debris will prevent the valve disc or poppet from moving at all.\nA need therefore exists for a simple, inexpensive yet reliable irrigation control valve which allows water to flow to an operational sprinkler, but restricts the water flow through a broken riser. Ideally, such a valve should allow detection of the broken riser, even when the flow is essentially shut off."}
{"text": "1. Field of Invention\nThis invention relates to golf and more specifically, to three-dimensional (3-D) mapping, photogrammetry, and light detection and ranging (LIDAR) techniques used to produce the software and data required for implementation into a series of golf course navigation software and physical hardware devices that use satellite navigation and geospatial database systems for assisting golf players, fans, and viewers in course navigation, visualization, management, and administration.\n2. Description of Related Art\nGolf is a game of accuracy, precision and distance. Without an expert caddie along, there's a lot of analysis and sometimes guesswork involved in determining shot distance to the hole, and various course targets such as greens or landing areas with water hazards in play. A golfer's club selection directly depends on the perceived distance of the desired shot.\nGlobal Navigation Satellite Systems (GNSS) is a standard generic term for satellite navigation systems that provide autonomous geo-spatial positioning with global coverage. GNSS allows small electronic receivers to determine their location (longitude, latitude, and altitude) to within several meters using time signals transmitted along a line-of-sight by radio from satellites. Receivers calculate the precise time as well as position, which can be used as a reference for distance measurements. Within the United States, the NAVSTAR global positioning system (GPS) is the only fully operational GNSS. GPS is currently the world's most utilized satellite navigation system. GPS technology has made its way into devices designed to be used as a navigation aid by a player on a golf course. The primary goal of a golf course navigation system is to provide players with distances to all kind of features on the course and especially to the green they are targeting. Beyond basic distance-to-the-hole data, many standard, golf GPS devices tell you the length of your last shot; distance to, and location of targets and hazards; and distance to the front, middle, and back of the green.\nConventional golf GPS devices suffer from a number of drawbacks. First, the geospatial data used (satellite imagery and differential GPS collection systems) in the production in commercially available golf GPS devices, such as those produced by GPS Industries, Sky Golf's Sky Caddie, and Callaway's uPro is very low resolution and is no better than five (5) to thirty (30) yards or meters. This degree of error is critical as a player's full shot distance typically varies by 10 yards or more for each golf club. For example, a player may hit an 8-iron 150 yards and a 7-iron 160 yards. If the golf GPS device indicates that the remaining distance to a pin is 150 yards when it is really 160 yards, then a player hitting an 8-iron instead of a 7-iron will come up well short—an unnecessary error that may result in the player's ball coming to rest in a hazard rather than on the green. This type of inherent error is not well suited for a game considered by many to be a game of inches.\nThe accuracy of conventional GPS devices is further complicated as some courses are better fit to use the capabilities of GPS than others. For example, in order to obtain a location reading, the GPS receiver must have a non-obstructed view of the sky. Accordingly, courses with lots of trees and other types of natural obstructions are less suitable for golf GPS devices. Objects such a mountains or buildings between the satellite and the receiver can produce significant error, sometimes up to 30 meters.\nThe graphics displayed on conventional GPS devices are rudimentary at best. Some devices are very basic as they merely display distance data without illustrating respective hole layout, e.g., a superimposed graphic displaying the overall geometry of the fairway, green, and related hazards. Other devices may show a hole's layout, but only through an artist's rendering in two-dimensions using computer generated imagery (CGI). These CGI renderings are then referenced geospatially with latitude and longitude values. For example, the artist will attribute the CGI created bunker with a lat/long value. When this data is loaded to a device with a GPS receiver, the GPS receiver will try to record a true location relating it to the pre-rendered imagery. Recently, airborne imagery of an “oblique” hole layout has been implemented in a golf GPS device provided by Callaway. However, this oblique imagery and perspective view was acquired by Pictometry's Oblique camera system and is not considered a mapping or calibrated camera system usable for mapping or civil engineering. Moreover, all of the images used in these systems are acquired using non-stereo specifications and are considered only two-dimensional in nature, i.e., the hole layout is only displayed as a flat two-dimensional image. No stereo compiled topographic or 3-D terrain data (x,y,z) are provided with respect to changes in topographic elevation, height, and related information such as slope. For a game that depends on where and the way a golf ball lands and bounces, conventional golf GPS systems are woefully inadequate. For example, any device that does not calculate the change in hypotenuse or change in elevation (topography) between point A and point B, is programmatically assuming that a golf course is flat and that a constant horizontal grid of X and Y values are sufficient. In fact, this is not the case as golf courses are not flat and a “true” golf course navigation device must calculate for the change in hypotenuse or change in elevation, topography or slope with respect to a calculated distance between point A and point B."}
{"text": "This invention relates broadly to the chemical modification of wool by reacting it simultaneously with a halogenated acid anhydride and an isocyanate. In particular, the invention concerns and has as its prime object the provision of processes wherein the reaction of wool with the above reagents is conducted in the presence of cresol, whereby to facilitate and promote the reaction. The unqualified term \"cresol\" used herein includes o-cresol, m-cresol, p-cresol, or any mixture of these isomers. Further objects and advantages of the invention will be apparent from the following description wherein parts and percentages are by weight, unless otherwise specified.\nAlthough wool is a very useful fiber, it is often desirable to improve its properties for particular applications by chemically modifying it. Previously, wool has been shrinkproofed with isocyanates, but the product is not flameproof. Furthermore, wool has been flameproofed with halogenated acid anhydrides, but the product does not meet consumer requirements for shrinkproofing.\nWool which is both shrinkproof and flameproof has also been prepared. However, a two-step process is required-- the wool is first shrinkproofed and then flameproofed. The two-step process has the disadvantage of being less efficient and economical than a one-step process."}
{"text": "In the field of magnetic recording using a head and a medium, further performance improvement of thin film magnetic heads and magnetic recording media is in demand in association with the high recording density of magnetic disc devices. As a thin film magnetic head, at present, a composite type thin film magnetic head made of a structure where a magnetoresistant (MR) element for reading and an electromagnetic transducer element for writing are laminated is widely used.\nThe magnetic recording medium is a discontinuous medium in which magnetic grains are aggregated, and each magnetic grain has a single magnetic domain structure. In the magnetic recording medium, one recording bit is configured by a plurality of magnetic grains. Consequently, in order to enhance the recording density, asperities at the border between adjacent recording bits need to be reduced by decreasing the size of the magnetic grains. However, if the magnetic grains are reduced in size, there is the problem that the thermal stability of magnetization of the magnetic grains is reduced in association with a decrease in the volume of the magnetic grain.\nAs a countermeasure against this problem, an increase of magneto anisotropy energy Ku of magnetic grains may be considered; however, the increase of Ku results in an increase in an anisotropy field (coercive force) of a magnetic recording medium. On the other hand, the upper limit of the recording magnetic field intensity of the thin film magnetic head is substantially determined by saturation magnetic flux density of a soft magnetic material configuring a magnetic core within the head. Consequently, if the anisotropy field of the magnetic recording medium exceeds the acceptable value determined by the upper limit of the recording magnetic field intensity, it becomes impossible to write to the magnetic recording medium. Currently, as a method to solve such a thermal stability problem, a so-called thermally assisted magnetic recording method is proposed in which a recording magnetic field is applied to record information under a state where the anisotropy field is reduced by heating the magnetic recording medium while a magnetic recording medium made from a magnetic material with large Ku is used.\nIn this thermally assisted magnetic recording method, a method using a near-field light probe, or so-called plasmon-generator, comprising a metal piece that generates NF light from plasmon excited by laser light is generally known, and as a magnetic head including such a plasmon-generator, a magnetic head including a magnetic pole, a waveguide and a plasmon-generator facing the waveguide is proposed.\nIn a thermally-assisted magnetic recording head, light propagated through the waveguide is coupled with a plasmon-generator in a surface plasmon mode, and excites the surface plasmon. Propagation of such surface plasmon in the plasmon-generator causes the generation of near-field light at the near-field light generating portion positioned at the end portion of the plasmon-generator at the medium opposed surface side. Then, a magnetic recording medium is heated by irradiating the magnetic recording medium with the near-field light generated at the near-field light generating portion of the plasmon-generator, and information is recorded by applying a magnetic field under a state where the anisotropy field of the magnetic recording field is decreased.\nIn a thermally-assisted magnetic disk device including such a thermally-assisted magnetic recording head, the distance between the thermally-assisted magnetic recording head and the magnetic recording medium, i.e., the flying height of the thermally-assisted magnetic recording head relative to the magnetic recording medium, is an important parameter to satisfy demands, such as reduction of the bit error rate (BER) or high recording density.\nRecently, in order to respond to demands, such as high recording density, a technology is proposed in which a medium opposed surface of the thermally-assisted magnetic recording head toward the magnetic recording medium side is protruded by heat-expanding the surface due to the heat generation of a heater and the like, and the flying height of the thermally-assisted magnetic recording head relative to the magnetic recording medium is reduced (see for example, JP 2010-79978). In the patent literature (JP 2010-79978), it is disclosed that the protrusion amount of the medium opposed surface is controlled by controlling the power supply to the heater, and the optimum flying height (the target flying height) can be realized. However, in a thermally-assisted magnetic disk device, in addition to the protrusion of the medium opposed surface due to supply of power to the heater and the application of writing current, since the medium opposed surface also protrudes due to the generation of near-field light, it is actually difficult to reduce the target flying height to approximately several nm (approximately 2 nm).\nIn the technology disclosed in the patent literature, the supply of power to the heater that can accomplish the target flying height can be set by a touch down test of the thermally-assisted magnetic recording head. In other words, the medium opposed surface gradually protrudes while the supply of power to the heater is increased, contact (touch down) between the thermally-assisted magnetic recording head and the magnetic recording medium is detected, and the supply of power to the heater that can accomplish the target flying height is determined.\nIn the meantime, in a thermally-assisted magnetic recording head, when a signal is recorded, the near-field light generated at the plasmon-generator is irradiated to the magnetic recording medium. In association with the generation of the near-field light, the vicinity of the medium opposed surface of the thermally-assisted magnetic recording head is heated, and the medium opposed surface protrudes toward the magnetic recording medium side. Therefore, taking into consideration the protrusion of the medium opposed surface of the thermally-assisted magnetic recording head due to the near-field light, it is necessary to set a supply of power to the heater that can accomplish the target flying height."}
{"text": "Microarray technology has been widely used for genomics and proteomics research as well as for drug screening. Currently, the spot size in most microarrays is larger than one micron. The use of nanometric biomolecular arrays, with smaller spot sizes, will enable high-throughput screening of biomolecules—eventually at the single molecule level. Also, nanometric arrays permitting precise control over the position and orientation of individual molecules will become a powerful tool for studying multi-valent and/or multi-component molecular interactions in biological systems. Toward these ends, protein arrays with feature sizes smaller than 100 nm have been fabricated, mostly using dip-pen nanolithography and nanografting. See Lee et al., Science 2002, vol. 295, p. 1702; Wilson et al., Proc. Natl. Acadi Sci. USA 2001, vol. 98, p. 13660; Liu et al., Proc. Natl. Acad. Sci. USA 2002, vol. 99, p. 5165; Pavlovic et al., Nano Lett. 2003, vol. 3, p. 779; and Krämer et al., Chem. Rev. 2003, vol. 103, p. 4367.\nBiological microelectromechanical systems (bioMEMS) are of tremendous interest for their potential applications in microscale, high throughput biosensing and medical devices (Shawgo et al., J. Curr. Opin. Solid State Mater. Sci. 2002, v. 6, p. 329). Using silicon as a substrate for the preparation of such devices is particularly attractive, since the extensive micro-fabrication techniques developed by the microelectronic industries can be used to fabricate and integrate various micro-components into the devices. For reducing biofouling, considerable research has been directed to the modification of substrate surfaces with stable and ultrathin films of poly(ethylene glycol) (PEG) or oligo(ethylene glycol) (OEG) (Prime et al., Science 1991, vol. 252, p. 1164). Since many of the ultimate applications for bio-devices require moderate-term (e.g., a few hours to several days) exposure to biological media (e.g., buffer of pH 7.4 at 37° C.), stability of the bio-compatible coatings on the devices under these conditions is highly desirable. All of the OEG/PEG-terminated films on silicon substrates reported by others are bound onto the silicon surfaces via Si—O bonds that are prone to hydrolysis (Calistri-Yeh et al., Langmuir 1996, v. 12, p. 2747), thereby limiting their stability under physiological conditions (Sharma et al., Langmuir 2004, v. 20, p. 348).\nAs described in commonly assigned, co-pending U.S. patent application Ser. No. 10/742,047, olig(ethylene glycol) (OEG) terminated alkenes were grafted onto hydrogen-terminated silicon surfaces through hydrosilylation (as developed by Linford and Chidsey, see Linford et al., J. Am. Chem. Soc. 1993, v. 115, p. 12631; Buriak, Chem. Rev. 2002, v. 102, p. 1271) forming robust Si—C bonds with the silicon surfaces. It was shown that the alkyl monolayers grown by this method were stable in boiling organic solvents, water, and acids, as well as slightly basic solutions (Linford et al., J. Am. Chem. Soc. 1995, v. 117, p. 3145). A method describing the modification of hydrogen-terminated silicon surfaces, including a silicon atomic force microscopy (AFM) cantilever tip, with OEG-terminated alkenes via either thermally- or photo-induced hydrosilylation is also found in commonly assigned, co-pending U.S. patent application Ser. No. 10/742,047. See also Yam et al., J. Am. Chem. Soc. 2003, v. 125, p. 7498; Yam et al., Chem. Commun., 2004, p. 2510). The efficiency with which such OEG-terminated films resist protein adsorption depends on many factors including the number of ethylene glycol (EG) units and the packing density of the films that is determined by the underlying substrate surface and the deposition methods. For example, OEG-terminated thiolate self-assembled monolayers (SAMs) on gold (111) surfaces are protein resistant, but those on silver (111) surfaces are not (Herrwerth et al., J. Am. Chem. Soc. 2003, vol. 125, p. 9359). The latter was attributed to the high packing density and structural ordering of the SAMs. Research has demonstrated that films grown on Si (111) surfaces had a density similar to that of the corresponding thiolate SAMs on gold (111) surfaces, and similarly reduced the adsorption of fibrinogen to 1% monolayer or less (Cai et al., to be submitted).\nAs a result of the foregoing, a nanometric biomolecular array comprising a stable, patternable monolayer surface, and an efficient method of making such an array, would be very beneficial."}
{"text": "The present invention relates to a presentation device for carrying out a product presentation. The present invention also relates to a method for presenting the product at a presentation location at which the customer is present.\nBefore a customer purchases a particular product, for example, a motor vehicle, he/she would generally like to obtain an impression of the article to be purchased. For this purpose, the customer may visit the vendor's premises. However, in markets covering a wide area, for example, in rural areas, the business premises of the commercial establishment and the domicile of the customer, for example, his/her place of residence, may be separated by great distances. This makes it unappealing for the customer to visit the business premises. Internet commerce allows customers to obtain an impression of the goods for sale conveniently from home. However, this impression is generally not highly realistic, and the photographs of the products depicted on the Internet pages often do not address the particular questions that a customer may have. For example, it is not possible to estimate how large a particular object is, based on a photograph of the product. For example, if the customer is interested in a motor vehicle and would like to know how spacious the rear seat is, he/she is not able to assess this based on the photographs.\nDE 10 2005 048 926 A1 describes a method for simulating the appearance of items of furniture at their place of destination. In this method, a customer at a home computer is able to create virtual models of both the items of furniture and the place of destination. The customer generates the model of the place of destination based on two-dimensional images, for example, photos of his/her living room. The virtual models of the item of furniture and the place of destination may then be combined with the aid of the computer.\nUS 2010/0289817 A1 also describes a method for illustrating a virtual object in real surroundings. In this method as well, the real surroundings are imaged by photographs, and a virtual model of the surroundings is formed from the photographs. The virtual objects may then be inserted into the model, overlaps of real objects imaged in the photographs with the virtual object being taken into account, so that the virtual object in the model may also be situated behind the real objects.\nDE 102 30 911 A1 describes a sales portal for the Internet, via which a customer may obtain product data which describe the product three-dimensionally, so that the user is able to superimpose images of the product into his/her field of vision by so-called augmented-reality glasses, thereby obtaining the impression that the product is in front of the user in his/her familiar living environment. However, to generate the three-dimensional virtual objects, the user needs a relatively high-powered computer as well as the expensive augmented-reality data glasses. Therefore, this purchasing portal is usable only by people who have both hardware components, the high-powered personal computer and the expensive data glasses."}
{"text": "There are many characteristics of tissue products such as bath and facial tissue that must be considered in producing a final product having desirable attributes for the product's intended purpose. Improved softness of the product has long been one major objective, particularly for the success of premium products. In general, the components of “softness” include stiffness and bulk (density), with lower stiffness and higher bulk (lower density) generally improving perceived softness. While enhanced softness is generally desirable for all types of tissue products, it has been especially challenging to achieve.\nThe throughdrying process is well known in the art and is used extensively in the manufacture of all types of tissue products having desirable softness characteristics. Throughdrying provides a relatively noncompressive method of removing water from a web by passing hot air through the web until it is dry. More specifically, a wet-laid web is transferred from the forming fabric to a coarse, highly permeable throughdrying fabric and retained on the throughdrying fabric until dry. The web is typically conveyed on the fabric through passage of one or more throughair dryers (“TAD”) where the web is dried to its final dryness state. The resulting dried web is typically softer and bulkier than a conventionally dried uncreped sheet because fewer bonds are formed in the web and the web is less compressed. Such throughdrying processes are described in detail, for example, in U.S. Pat. Nos. 6,149,767 and 6,331,230 B1.\nHowever, from a manufacturing perspective, the throughdrying process is relatively energy intensive and therefore expensive compared to conventional wet pressing processes. The energy consumption issue is compounded by the fact that a significant amount of heat is vented to atmosphere during the process and, thus, wasted.\nThe present invention relates to an improvement in throughdrying processes that will significantly enhance the energy efficiency of such processes."}
{"text": "It is well known to establish a vehicle occupant classification such as, for example, by measuring the weight of an occupant. Vehicle occupant classification information identifies the type of occupant seated within a vehicle and generally includes the following three categories: adult, child, or none. This information may be useful, for example, in determining whether or not to deploy an airbag. As an example, it may be desirable to deploy an airbag under certain circumstances if the vehicle occupant in a particular seat is an adult, but the airbag may not be deployed if the vehicle occupant is a child or if the particular seat is empty."}
{"text": "This application is a continuation of pending U.S. patent application Ser. No. 13/184,619, filed on Jul. 18, 2011 and entitled “Methods And Apparatus For Vacuum Ultraviolet (VUV) Or Shorter Wavelength Circular Dichroism Spectroscopy” which claims priority to Provisional Patent Application No. 61/400,153 filed Jul. 22, 2010; the disclosures of which are expressly incorporated herein by reference.\nThe present disclosure relates to the field of polarization spectroscopy. More specifically, it provides a means by which circular dichroism (CD) measurements may be performed in the vacuum ultraviolet (VUV). In one embodiment a highly efficient laboratory-scale VUV CD spectrometer is provided. As used herein, vacuum ultraviolet (VUV) light includes, generally, wavelengths of light that are about 190 nm and less wavelengths.\nWhen light passes through a solution of optically active (e.g. chiral) molecules, left- and right-circularly polarized components traverse the solution with different speeds and are absorbed by the solution to different extents. The differing propagation speeds lead to the effect of optical rotation (OR), which can be probed by measuring the rotation of the polarization plane via optical rotatory dispersion (ORD). Near absorption bands, circular dichroism (CD) spectroscopy measures the differential absorption of left- and right-circularly polarized light by the substance. Both methods have been successfully used to characterize properties of solutions containing optically active molecules. In addition to being sensitive to inherent chirality, ORD and CD are also sensitive to different conformations of complex molecules. As a result, one of the most prolific applications of CD spectroscopy has been the study of the secondary structure of soluble proteins.\nOptical rotation and circular dichroism have the same underlying cause, and are related to each other via Kramers-Kronig integration. Defining circular birefringence (CB) as twice the optical rotation, a complex general circular retardation can be defined by:C=CB−iCD,  Eqn. 1where CB and CD are functions of wavelength. The form of Eqn. 1 and the Kramers-Kronig relationship it obeys are familiar from the (isotropic) optical dispersion of materials, n−ik, where n is the material's refractive index and k is the extinction coefficient. That n and k obey the Kramers-Kronig relation implies that the wavelength dependence of the index of refraction is determined by the material's absorption, and vice versa. Similarly, that CB and CD are related via Kramers-Kronig integration implies that complete knowledge of one over all wavelengths or energies determines the other.\nPractical experiments, however, occur over finite wavelength ranges and ORD and CD spectroscopy have different advantages, depending on the wavelength range explored. While the ORD spectrum can be determined far from the absorption bands responsible for its effect, the CD spectrum directly probes the bands, and as such, is considered more sensitive and spectrally “compact”. In ORD spectroscopy the information associated with a given absorption band is typically spread out over a large energy range. As a result, the ORD spectrum in a given wavelength range contains overlapping contributions from multiple absorption bands. In contrast, with CD spectroscopy the information due to a single absorption band is localized within a smaller energy range, resulting in much less overlap between measured features. If the absorption region is accessible experimentally, CD spectroscopy is typically the preferred technique since it more directly probes the underlying “cause” of the optical activity of the substance being measured.\nIn the case of soluble proteins, the secondary structure of the molecules results in the optical activity of the solution, and different types of secondary structure, such as alpha-helix, beta sheet, polyproline-II helix, etc. give rise to distinct features in CD spectra. The electronic bands responsible for protein CD largely reside in the ultraviolet. The alpha-helical structures generally involve bands centered at ˜222, 208, 192, 175, 160, and 140 nm. Meanwhile, beta sheet structures have weaker bands at ˜215, 198, 175, and 168 nm. Variations in beta sheet geometry result in further modifications to CD spectra. The polyprolene-II helix has observed bands near 226 nm and 206 nm, and gives rise to spectra similar to those of proteins previously characterized as “un-ordered”. In all cases, other bands may exist but have yet to be observed and/or identified.\nThe ultraviolet wavelength region may be considered as consisting of two distinct segments: the near-UV region from ˜190-400 nm, and the “far-UV” or vacuum ultraviolet (VUV) region below 190 nm. Conventional CD spectrometers are limited in operation down to about 190 nm. The primary motivation to extend CD studies into the VUV lies in the existence of additional absorption bands present at shorter wavelengths. As a consequence of these supplementary features, VUV CD spectra inherently posses increased information content relative to their longer wavelength counterparts. While traditional CD instruments are often limited to determining the amount of alpha-helix structure present in a solution, VUV CD systems are capable of extracting numerous secondary component fractions. In addition, these powerful systems can also provide insight into conformational changes, such as fold state, independent of secondary structure. It is important to note that this information enhancement is not simply due to a “more data is better” argument; the improved capabilities are a direct consequence of the additional absorption bands present in the VUV.\nOptical studies in the VUV are difficult to conduct due to the intrinsic absorption of most materials in this region. This phenomenon precludes the simple extension of, or modification to, traditional longer wavelength optical instrumentation to facilitate operation at these energies. To achieve efficient optical performance in the VUV, an instrument must be explicitly designed to do so. Specifically, conventional optical systems are designed to operate in atmospheric conditions and typically lack, among other things, the controlled environment required for operation at these shorter wavelengths. VUV radiation is strongly absorbed by both oxygen and moisture; hence, these species must be maintained at sufficiently low levels in order to permit transmission of VUV photons through instrument optical paths. Attempts to reach shorter wavelengths by simply purging with non-absorbing gases generally yield poor results. Furthermore, transmissive optical components that are otherwise suitable for near-UV or visible wavelength operation, routinely absorb strongly in the VUV. Consequently, reflective elements must instead be employed, greatly restricting design options. As a result, it is comparatively difficult to achieve high optical throughput in VUV optical instrumentation.\nWith the exception of early work by Johnson (Johnson W. C. (1971). “A circular dichroism spectrometer for the vacuum ultraviolet”. Rev. Sci. Instrum. 42(9): 1283-1286), progress towards development of a dedicated, highly efficient VUV-CD instrument has been limited. Today's commercial bench top systems are designed to operate at near-UV and longer wavelengths. Several of these systems offer simplistic purge functionality in an effort to extend capabilities into the VUV. Practically however, poor data quality restricts these instruments to operation at wavelengths above approximately 185 nm.\nThe most capable VUV CD systems in existence are those integrated into synchrotron radiation (SR) beam lines. The advent of such instruments in the 1980's and 1990's brought about tremendous enhancements in both CD data quality and information content. These improvements however, were largely the result of the remarkable intensity of SR sources, rather than fundamental advancements in instrument designs. As a consequence of these developments, interest in VUV-CD spectroscopy has grown considerably, resulting in the commissioning of several new beam lines and the identification of a myriad of applications for this technology. SR-CD systems are powerful, but the disadvantages are obvious: synchrotron facilities are huge and enormously expensive, making accessibility a severe limitation.\nIt follows that there would be great benefit in the development of a highly efficient, high throughput laboratory-scale VUV CD spectrometer, which does not require synchrotron radiation. Such an instrument would render high-throughput structural investigations of proteins widely available, thus creating new opportunities to accelerate discoveries in structural proteomics."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nRiding mowers typically include some form of steering device. These steering devices can include hydraulic cylinders and electric rotary motors for wheel turning operation. However, these steering devices can include complicated geometry including a large number of parts. Further, use of hydraulic cylinder devices can result in oil spillage, which can damage turf."}
{"text": "The present invention is generally related to radiotelephones, and more particularly to a cellular portable telephone which accommodates a power source change and corresponding operating mode change during a radiotelephone call.\nCellular portable telephones currently have batteries with capacity for one-half hour of continuous operation. As a result most users typically carry one or more spare batteries so that they can switch to another battery when the portable telephone indicates that the battery voltage is low. However, when the battery is removed during a cellular telephone call, the call is terminated and the user must redial the call. Interruption of the power source and termination of a cellular telephone call in process also occurs when the user plugs the cellular portable telephone into a vehicle's battery by way of a vehicular adaptor installed therein, such as a battery eliminator connected to the vehicle's cigarette lighter. For the foregoing reasons, there is a need for a cellular portable telephone which accommodates a power source change and corresponding operating mode change during a cellular telephone call."}
{"text": "The present invention relates to a screen for treating fibrous suspensions, such as pulps, of the wood processing industry. Especially it relates to the construction of a rotor element for the screen.\nPressure screens are essential devices in the production of pulp and paper. They remove from the pulp suspension mainly impurities, over-sized pieces of wood and fiber bundles as well as other undesired substances. The screen can also fractionate fibers according to their length for improving the properties of the pulp. The precise function of the screen is dependent on the location in the process where it is used. In the screening process the water suspension of the pulp fibers is typically pumped into a cylindrical chamber, wherein the suspension is brought to contact with the screen surface and a rotor moving at high velocity. The rotational velocity of the rotor pushes the fibrous material into movement, whereby part of it is passed as accept through apertures in the screen surface. The high-speed rotor applies positive and negative impact pulses to the suspension. The positive impact pulses push the fibers through the apertures in the screen and may fractionate the fibers. The negative impact pulses provide for a regular flush-back of the apertures in the screen surface so that the fibers do not plug the apertures.\nThe pulp suspension consists of millions of elastic fibers that easily attach to each other forming so-called fiber flocks. Even at a low consistency such as 0.01% the fibers form unstable flocs. In a typical screening consistency, 1-3% the fibers form stable flocks and fiber networks hamper the screening. The fibers and undesired solid matter are periodically removed from the net in order to enable the screening the remaining fibers from the flocks and fiber networks into reject and accept fibers. When the pulp consistency increases, the force required for decomposing the fiber network increases intensively and finally a process limit is reached, where the apertures in the screen surface or the reject line is clogged. A large number of various rotor solutions has been developed with the aim of ensuring a continuous screening operation.\nIn principle, the rotors can be divided into two basic groups, open and closed rotors. Both are being used and their purpose is, as known, to keep the screening surface clean, i.e. to prevent the formation of a fiber mat on the screening surface. The first group is characterized in that the interior of the screen drum is provided with a rotary shaft or a rotor, whereto blades are attached by means of arms. An example of this kind is the rotor solution according to U.S. Pat. No. 4,193,865, where the rotor is arranged rotatably inside a cylindrical stationary screen drum, said rotor comprising blades located in the vicinity of the screen drum surface, which blades in the construction according to said patent form an angle with the drum axis i.e. the blades extend obliquely from one end of the screen drum to another. When moving, the blades impact pressure pulses on the screen surface, which pulses open the surface apertures. There are also solutions, in which the blades have been located on both sides of the screen drum. In that case, the suspension to be treated is fed to the inside or to the outside of the drum and the accept is, respectively, discharged from the outside or inside of the drum.\nIn stationary rotors the rotor is an essentially closed cylindrical piece, the surface of which is provided with pulsation members, for instance almost hemispherical protrusions, so-called bulges. In this kind of an apparatus the pulp is fed into a treatment space located between the rotor cylinder and the screen drum outside thereof, whereby the purpose of the rotor protrusions, e.g., the bulges, is both to press the pulp against the screen drum and by means of its trailing edge to withdraw the fiber mat off the screen drum apertures. The bulges can be replaced by other kinds of protrusions.\nA solution widely used in the market is a represented by a method according to FI patent 77279 (U.S. Pat. No. 5,000,842) and the solution developed for the implementation thereof. The method according to said patent is characterized in that the fiber suspension is subjected to axial forces with varying intensity and effective direction, the direction and intensity of which are determined based on the mutual axial positioning of the point of application and the countersurface of the screen drum and by means of which the axial velocity profile of the fiber suspension is changed while maintaining the flow direction continuously towards the discharge end. Preferably the surface of the rotor is divided into four zones: feed, feed and mixing, mixing, and efficient mixing. The rotor surface is typically provided with 10-40 protrusions, the shape of which varies according to the zone i.e. the axial part of the rotor that they are located on. The protrusions on the housing surface of the rotor are mainly formed of front surfaces facing the flow, preferably surfaces parallel to the housing surface and back surfaces that descend towards the housing surface of the rotor. The housing surface of the rotor is provided with protrusions of several different forms, which have been arranged onto the rotor housing so that two or more circumferential zones are formed separated from each other in the axial direction of the rotor, such as e.g. 4 zones. At least part of the front surfaces of the protrusions forms an angle with the axial direction. The front surface of the protrusions can be divided into two parts that form with the axial direction angles of different size. The variation interval of the angles is −45°-+45° compared to the axial direction. However, the functioning principle of the protrusions is the same as in other corresponding devices. The abrupt front surface imparts a strong pressure shock to the fiber mat on the screen drum, whereby the accept is pressed through the apertures of the drum. The sloping back surface of the protrusion withdraws some water back to the screening zone and thus releases from the grooves and apertures major particles and fiber flocks thus cleaning the screen drum.\nU.S. Pat. No. 5,192,438 describes a rotor which provides high intensity axial shear stress in addition to high positive pulses and negative pulses. The rotor has a contoured surface including a plurality of protrusions. A protrusion has a front plane, an upper plane, an inclined plane and edge surfaces, which may converge. The trailing surface of the protrusion is abrupt.\nSo, in prior known solutions the functional prerequisite of pressure screens starts from the presumption that the rotor element is to develop an adequate pressure impulse on the interface to make the fiber particles flow through the screening surface and that the rotor element is to create by its trailing edge a negative pressure impulse to generate a turbulence that cleans the apertures clogged by the previous positive impulse. It has also been generally presented in the field that a negative impulse withdraws liquid back towards the feeding space preventing excess thickening of the fiber suspension in the feeding space and in its part cleaning the apertures of the screening surface. For enabling to create these conditions, the rotor must have an adequate rotational speed, which is, however, limited by energy consumption and mechanical durability of the screen, a typical speed for a rotor described in FI-patent 77279 (U.S. Pat. No. 5,000,842) is 24 m/s.\nIn the present industrially used pressure screen applications the rotor solutions have enabled to reach the maximum feed consistency level of pulp. The consistency level is almost the same for different rotor types, for instance for softwood (SW)-pulp approximately 2-3%. Thus, there is a need in the field to develop a screen rotor that will allow higher feed consistencies."}
{"text": "The various embodiments described herein provide a system and method for dynamically rendering a 3D representation of data in a dimension different than an initial 2D base layer representing related space and aerospace datasets, although the techniques described herein will find broader occasions in a variety of data representation situations.\nIn an embodiment, the user selects desired parameters from a datastore to be displayed in the 2D base layer. The user also selects one or more parameters to be represented in a third dimension from a datastore. The system has an inventory of display techniques through which certain data may be displayed.\nThe system also comprises rules for the display of specific data. For example, and without limitation, certain data may be appropriately displayed in one format but is not susceptible to being displayed in a second format. The processor and server have rules associated with the data type to allow the processor to present the user options for display of a specific data based upon the techniques by which that data may be displayed. Thus the user is presented only with choices for display that are relevant to the data type.\nBased on the type of data selected by the user, the processor determines which of a plurality of predetermined display formats would be appropriate and presents the options to the user. The user selects the desired display format from the list of choices and the processor generates the image.\nThe various display formats are selected from the set comprising 2D surface morphing, 2D plane rotation, 2D to 3D continuous morphing, 3D to 2D continuous morphing, and 2D/3D combined data representation. This list is by way of illustration only and is not meant as a limitation."}
{"text": "The quantitative determination of analytes in body fluids is of great importance in the diagnoses and maintenance of certain physiological abnormalities. For example, lactate, cholesterol and bilirubin should be monitored in certain individuals. In particular, determining glucose in body fluids is important to diabetic individuals who must frequently check the glucose level in their body fluids to regulate the glucose intake in their diets. While the remainder of the disclosure herein will be directed towards determining glucose, it is to be understood that the methods of this invention may be used for determining other analytes on selection of an appropriate enzyme.\nThe results of such tests can be used to determine what, if any, insulin or other medication needs to be administered. In one type of blood glucose testing system, sensors are used to test a fluid such as a sample of blood.\nA sensor contains biosensing or reagent material that will react with blood glucose. The testing end of the sensor is adapted to be placed into the fluid being tested, for example, blood that has accumulated on a person's finger after the finger has been pricked. The fluid is drawn into a capillary channel that extends in the sensor from the testing end to the reagent material by capillary action so that a sufficient amount of fluid to be tested is drawn into the sensor. The fluid then chemically reacts with the reagent material in the sensor resulting in an electrical signal indicative of the glucose level in the fluid being tested is supplied to contact areas located near the rear or contact end of the sensor.\nSuch a sensor is often sensitive to the effects of ambient humidity. One way to reduce or eliminate the effects of ambient humidity is to individually package each of the sensors with desiccant. Such a method has a drawback of requiring the unpacking of a strip before each use. Thus, it would be desirable to have a cartridge that would contain a plurality of test sensors that would not require unpacking each strip before using. Also, for the convenience and case of use, it would also be desirable to have a simple mechanism to feed the test sensors one at a time for testing by the user. This provides ease of use to normal users and is especially important for those users who may have some physical limitations."}
{"text": "1. Field of the Invention\nThis invention relates to a power semiconductor component having at least three layers, including first and second layers extending up to a first surface which has sectionalized contact planes resulting in a step-like structure for contacting the first and second layers, and to a process for its manufacture.\n2. Description of the Prior Art\nSuch a component as noted above is known, for example, from German Offenlegungsschrift No. 2,719,219. The semiconductor device described therein has four layers in a pnpn arrangement. The n-doped cathode layer 8 is subdivided into several sections, in order to allow contacting in two planes (FIG. 1). It has a surface doping of 5.times.10.sup.20 atoms cm.sup.3 and is 10 .mu.m thick. The p-base layer 6, located underneath, has a surface doping of 2.times.10.sup.18 atoms/cm.sup.3 and is 53 .mu.m thick. The width of the cathode sections is at least 320 .mu.m. The maximum reverse voltage of the pn junction between the cathode layer 8 and the p base 6 is 10 volt. The semiconductor device described is designed for a total current capacity of 400 A.\nA process for manufacturing the abovementioned semiconductor device is also indicated (FIG. 6). For this purpose, boron or gallium is first diffused into both sides of an n-doped silicon substrate 42. The doping of the n substrate is about 10.sup.13 -10.sup.16 atoms/cm.sup.3. The p layers which have been diffused in, the p base layer 46 and the anode 44 have a surface doping of about 2.times.10.sup.18 atoms/cm.sup.3. To form the cathode layer 48, phosphorus is diffused into the p base layer 46. This n layer 48 is photolithographically subdivided into several sections. Aluminium layers are then applied to the anode layer 44, to the cathode sections 48 and to the subdivided p base layers 46.\nThe known semiconductor device can be switched off by applying a negative current pulse to the p base layer 6 (gate). The power required for this purpose and supplied from the outside is considerable. The load current is constricted towards the middle of a cathode section by this current pulse and a part of the load current is thus drawn away to the gate. The current density of the remaining load current then becomes very high. If the switching-off step proceeds too slowly--for example in the case of unduly high switching-off amplification--or not homogeneously over the entire area of the semiconductor device, the latter is locally overloaded thermally and hence destroyed. The switching-off current via the gate produces a lateral voltage drop which should be smaller than the breakdown voltage of the pn junction between the cathode and the p base layer. However, the maximum reverse voltage of this pn junction is so low that an avalanche breakdown of this junction can occur during the switching-off step. The manufacturing process described is admittedly adequate for the structure having the dimensions of the known semiconductor device, but finer structures can hardly be produced satisfactorily this way.\nGerman Offenlegungsschrift No. 2,855,546 describes a field-controlled thyristor which has characteristics similar to those of a p-i-n diode. This thyristor consists of a thin, heavily p-doped anode zone 15, a thick, lightly n-doped semiconductor substrate 11, several deep p-doped grid zones 12 and thin, heavily n-doped cathode zones 14 arranged in between (FIG. 3). The grid zones 12 consist of channels having a width of about 12 .mu.m and a depth of 15 to 40 .mu.m. To produce these grid zones 12, channels are etched into a &lt;110&gt;-oriented silicon substrate with a mixture of potassium hydroxide and isopropanol in a ratio of about 3:1 (FIG. 1). For this purpose, a mask 13 of silicon dioxide is first formed. The channels are then epitactically refilled with heavily p-doped silicon. A thin layer of poly-silicon then additionally forms on the silicon oxide layer which is removed by etching. A mask of silicon oxide is then formed for the cathode zones 14 which are produced by diffusion of phosphorus. Finally, the grid zones 12 and the cathode zones 14 are provided with metallizations, silicon dioxide being left in between for electrical insulation.\nThe silicon with &lt;110&gt;-orientation, required for the process described, cannot be drawn without dislocations. For this reason, the silicon wafers can be obtained only in an elliptical form and must be sand-blasted to give a round form for further processing. The base material for the component described is therefore much more expensive than that for conventional components."}
{"text": "This invention is directed to a new and improved toy device for launching a compressible foam ball.\nIn the past, many different toy constructions for launching balls were proposed e.g. see U.S. Pat. Nos. 1,171,197; 2,725,869; 2,853,991; 3,120,387; 3,236,521; 3,301,246; 3,744,472 and 3,765,396.\nThe present invention is an improvement over such prior art based on its simple construction as well as the manner in which it functions to propell a soft compressible closed cell foam ball from the forward section of the lanucher."}
{"text": "1. Field of the Invention\nThe invention relates to semiconductor wafer processing equipment and, more particularly, the invention relates to a method and apparatus for selectively marking a semiconductor wafer with identifiable markings.\n2. Description of the Background Art\nIdentifying defects on the surface of a wafer, such as particulate contaminants and other surface irregularities, is extremely important to integrated circuit manufacturing processes. To eliminate defect sources, defects are identified and analyzed to determine the source of the defect. Thereafter, corrective action can be taken to reduce or eliminate the defect source.\nTypically, the defect identification process is accomplished in two steps. First, a laser scanner device scans a wafer with a laser and analyses the backscatter of the laser to locate defects on the surface of the wafer. One such laser scanner device is a Tencor SurfScan 6200, manufactured by Tencor Instruments. Second, each defect located by the scanner is analyzed to identify the root cause of the defect. The most commonly employed analysis tool is a high magnification imaging system such as a scanning electron microscope (SEM). An SEM is used to identify the defect and/or the source of the defect by inspecting the defect at high magnification. Additionally, the SEM may be accompanied by instrumentation for performing chemical analysis of the defect. Such instrumentation includes an energy dispersive x-ray (EDX) detector. Other tools may include instrumentation for performing an Auger analysis, an atomic force microscope (AFM), a tunneling electron microscope (TEM), an optical spectrometer and the like.\nSince the SEM as well as other analysis tools use a high magnification (on the order of 200 to 2000 times), rapidly positioning the SEM at the defect location can be time consuming. Although the laser device provides defect coordinates that are accurate to approximately 300 microns, a substantial amount of time can be spent manually searching, for example, a 300 by 300 micron region with the SEM for a defect having dimensions of a few tenths of micron.\nAdditionally, when using a SEM (or other optical analysis tool) to analyze a \"bare\" wafer, i.e., a wafer having no surface features, the SEM has difficulty focusing on the wafer surface. Without an accurate focus, finding a small defect on the surface is nearly impossible.\nFurthermore, some defects are identified by laser scanning, but cannot be seen with a SEM. As such, the SEM operator may search for a defect for a long period of time until realizing the defect cannot be seen with the SEM.\nSome defect analysis tools, e.g., an atomic force microscope (AFM), are not readily useful when a large search area is used. Generally, the field of view for an AFM is approximately 10 .mu.m.sup.2, and the search area is as large as 300 .mu.m.sup.2. Unfortunately, an AFM requires approximately five minutes to obtain a 10 .mu.m.sup.2 image. As such, 900 images are required to cover a 300 .mu.m.sup.2 search area, requiring 4500 minutes to complete the search.\nTherefore, a need exists in the art for a method and apparatus that selectively marks a wafer proximate a defect such that the defect can be rapidly identified and examined with a SEM or other defect analysis system."}
{"text": "1. Field of the Invention\nThis invention relates to a terpolymer with a high refractive index and the plastic lenses made of the terpolymer.\n2. Description of the Prior Art\nRecently lenses of organic glass have been widely used instead of lenses of inorganic glass, because organic glass is superior to inorganic glass in its lightness, safety due to its high impact resistance, workability, dyeability, etc. For example, lenses made of diethylene glycol bisallyl carbonate (hereinafter referred to CR39) have already been used widely. But, lenses of CR39 have such a disadvantage that the lenses must be thicker than those of inorganic glass, because of their relatively low refractive index such as 1.499. There is known a method for overcoming the disadvantage comprising copolymerizing CR39 with benzyl methacrylate so as to increase the refractive index. However, lenses obtained by the copolymerization are inferior to lenses of CR39 in their superficial hardness, so that the lenses must be more carefully handled.\nWe have now found that copolymerization of ternary monomers consisting of CR39, benzyl methacrylate, and triallylisocyanurate enables to produce lens with a high refractive index and a high superficial hardness.\nAn object of this invention is to provide a plastic lens with a high refractive index and a high superficial hardness."}
{"text": "In computer graphics, three-dimensional (3D) modeling involves generation of a representation of a 3D surface of an object. The representation may be referred to as a 3D object data model, and can be rendered or displayed as a two-dimensional image via 3D rendering or displayed as a 3D image. 3D object data models represent a 3D object using a collection of points in 3D space, connected by various geometric entities such as triangles, lines, curved surfaces, etc. Various techniques exist for generating 3D object data models utilizing point clouds and geometric shapes, for examples.\nBeing a collection of data, 3D models can be created by hand, algorithmically, or objects can be scanned, for example. As an example, an artist may manually generate a 3D image of an object that can be used as the 3D model. As another example, a given object may be scanned from a number of different angles, and the scanned images can be combined to generate the 3D image of the object. As still another example, an image of an object may be used to generate a point cloud that can be algorithmically processed to generate the 3D image.\n3D object data models may include solid models that define a volume of the object, or may include shell or boundary models that represent a surface (e.g. the boundary) of the object. Because an appearance of an object depends largely on an exterior of the object, boundary representations are common in computer graphics.\n3D models are used in a wide variety of fields, and may be displayed using a number of different types of interfaces. Example interfaces may provide functionality to enable interaction between a user and the 3D models."}
{"text": "The present invention is directed to the field of serving food portions, and is more specifically directed to a device for inserting a filling into a taco shell or the like. Taco shells, as found in restaurants or as packaged for sales to consumers in food stores, have a substantially uniform shape and size. However, the shape, which in profile is substantially semi-circular, does not lend itself easily to filling with presently available devices. For example, a teaspoon, tablespoon, or large serving spoon can be used to deposit the filling in the taco shell, but this generally results in the entire portion of the filling being deposited in the center, and none in the sides. If a spoon is used, additional portions of filling can be inserted on either side, but because of the semi-circular shape, the filling will tend to fall out at the sides. Likewise, if the spoon is used to more evenly distribute the filling across the bottom of the taco shell, there will be too much filling at the sides, which will tend to fall out, and not enough filling at the center.\nThere are many specialized devices for serving specific kinds of food, but I have found that none of these are suitable for filling taco shells. For example, a number of different devices are known for serving a portion of food such as ice cream or hamburger meat. A common type of device for serving ice cream includes a scoop, a handle extending rearwardly from the scoop, and a plunger with an ejector bar mounted for reciprocating motion in the handle. The shape of the scoop is dependent upon the shape of the ice cream portion which is desired. Thus, the scoop is often triangular, to produce a portion which will fit into an ice cream cone or a tall ice cream dish. In such devices, the end of the scoop which is inserted into the ice cream forms the point of the triangle, to make introduction of the scoop into the hard ice cream easier. Such devices are characterized by U.S. Pat. Nos. 1,688,595 and 1,798,490 to Parr and 2,178,648 to Rothenbush. In other devices, the scoop is rectangular in shape, to form a rectangular brick for filling an ice cream sandwich. Such devices are characterized by U.S. Pat. Nos. 1,352,755 to Levy and 1,529,782 to Gerstein. Although all of these devices are excellent for their stated purpose, that is, serving portions of ice cream, the sizes and shapes of the various scoops prevent turning them to use in filling a taco shell.\nSomewhat analogously, devices have been provided for slicing an opening in a sandwich bun and inserting into the opening the filling, for example a hamburger or a hot dog. Again, the shape of the scoop is dependent upon the shape of the sandwich bun, and the front is provided with a pointed blade for cutting the opening in the sandwich. Such a device is characterized by U.S. Pat. No. 2,784,682 to Clevenger. The size and shape of this type of device also renders it unsuitable for use in filling a taco shell.\nIn summary, no simple device exists which in shape and size is adapted to fill a taco shell with an amount of filling which corresponds to the shape of the taco shell. It is the solution of these problems to which the present invention is directed."}
{"text": "Field-effect transistors include source, drain and gate structures. A biasing voltage applied across gate and source terminals allows the flow of charge carriers, namely electrons or holes, between source and drain. High electron mobility transistors (HEMTs) are characterized by heterojunctions having different materials on either side of the junctions. Electrons remain close to the heterojunction and form two dimensional electron gas (2DEG). A bias applied to the gate modulates the number of electrons and controls the conductivity of the transistor.\nActive matrix devices such as displays (e.g. televisions, laptop monitors), imagers (e.g. x-ray imagers) and sensors typically use hydrogenated amorphous silicon (α-Si:H) and, in some applications, low-temperature poly-silicon (LTPS) thin-film transistor (TFT) backplanes on glass or, for flexible devices, clear plastic. However, for very high resolution applications (>1000 pixels per inch (ppi)), such as micro-displays or pico-projectors, the carrier mobility of a-Si:H (electron mobility of about 1 cm2/Vs) is too low to provide sufficient drive current at short TFT channel widths. For applications requiring high drive current, such as active matrix organic light emitting diode (AMOLED) displays, it is necessary to shrink the gate length and/or increase the gate width of a-Si:H transistors. This leads to increasing the processing cost of a-Si:H active matrix circuits due to the relatively small gate lengths as well as a significant trade-off in display resolution due to larger gate widths. LTPS is more expensive than a-Si:H, but capable of providing higher drive currents. The device-to-device variation of threshold voltage and mobility in LTPS transistors requires compensation circuitry that limits the resolution of the active matrix. Single crystalline silicon (c-Si) has been used as an alternative for very high resolution backplanes, but processing c-Si can require high temperatures not compatible with glass substrates currently used in manufacturing a-Si:H or LTPS devices or clear plastic substrates that may be used.\nSome existing displays have pixel densities of about 100 PPI (pixels per inch), each pixel including three RGB sub-pixels. Pixel dimensions of such devices may be about one hundred microns (100 μm). Such displays further include organic light emitting diodes (OLEDs) requiring a drive current of about 300 nA for a 100 μm pixel. Amorphous hydrogenated silicon thin film transistors (TFTs) having standard SiNx gate dielectrics are employed in conjunction with the OLEDs. Using such TFTs, resolutions greater than 150 PPI are difficult."}
{"text": "Array antennas are used for example in satellites of a satellite system for worldwide navigation (GNSS; GNSS=Global Navigation Satellite System), such as the currently constructed European satellite navigation system, better known under the name “Galileo”, in satellite-based imaging radar systems employing Synthetic-Aperture-Radar (SAR) sensors or in Telecommunication satellites. In order to control the antenna excitation, the characteristics of each branch of the distribution network and array radiator elements of the array antenna must be known. Typical parameters describing the characteristics are the propagation time, the phase and the magnitude of a signal propagating through each branch of the array. The characteristics may be used to calibrate the different paths or branches of an array antenna, thus enabling an improved control of the excitation. Usually, an estimation or determination of the characteristics and a calibration cannot be performed during normal operation of an antenna, since this requires emitting special test or calibration signals, which would interrupt the normal operation. This aspect is particular important in case of a GNSS, where the system requirements of operational continuity are usually mandatory."}
{"text": "This invention relates to a device for providing traction to the powered wheels of a vehicle. The device is useful for extricating a vehicle which has become stuck on a slippery surface, such as ice, snow, or mud, as well as for assisting a vehicle which is to be driven on such surfaces.\nPreviously, there have been many devices which were adapted to be attached to the powered wheels of a vehicle to increase its traction. Such devices included conventional chains and many cleated devices. One of the problems with such past devices was that they usually included means for pulling the cleats tightly against the tire. This tight arrangement caused the cleats to constantly rub against and to ultimately destroy the tire.\nAnother of the problems with such past devices was that the cleats were often fixed to or urged against the cleat holder. This tight arrangement also caused the cleats to constantly rub against and to ultimately destroy the tire.\nAnother of the problems with such past devices was that the devices were difficult or impossible to mount onto the tires unless the tires were lifted above the driving surface. Obviously, such lifting is not possible when the vehicle is stuck in a rut.\nAnother of the problems with many such past devices was they they required two cleat holders, one on the outside and one on the inside of the tire. This caused the person installing the device to have to reach around to the inner side of the tire (under the vehicle) and to have to work on apparatus elements that he could not see.\nThe object of this invention is to overcome all of the foregoing problems by providing a tire traction device that is maintained loose on the tire and loose on the cleat holder; that does not destroy the tire; that does not require the tire to be lifted off the ground; and that does not require a second cleat holder to be used on the inner side of the tire. It is another object of this device to achieve the foregoing object while simultaneously providing the vehicle with greatly increased traction at an inexpensive price."}
{"text": "The present invention relates to an X-ray CT (Computed Tomography) imaging method, subject center position detection method and X-ray CT apparatus, and more particularly to an X-ray CT imaging method, subject center position detection method and X-ray CT apparatus which can detect whether the center position of a subject is significantly offset from a scan center.\nIn X-ray CT imaging, a human operator adjusts the height of a top plate for laying a subject by visual measurement so that the center position of the subject coincides with a scan center (a center of rotation of an X-ray tube) with minimum offset.\nHowever, the adjustment by the operator\"\"s visual measurement sometimes makes the center position of the subject significantly offset from the scan center, which results in poor homogeneity in the CT value of a subject\"\"s image. Moreover, overrange may occur near the scan center, when the scan center falls outside the subject.\nIt is therefore an object of the present invention to provide an X-ray CT imaging method, subject center position detection method and X-ray CT apparatus which can detect whether the center position of a subject is significantly offset from a scan center.\nIn accordance with a first aspect of the present invention, there is provided an X-ray CT imaging method for performing a scan by an X-ray tube rotating around a subject, comprising the steps of: detecting a subject center position based on a projection profile acquired by imaging the subject; and giving a warning when a calculated distance between the subject center position and a scan center exceeds a predefined range.\nThe X-ray CT imaging method of the first aspect displays a warning on a display screen or gives a vocal prompt to a human operator to adjust the position when the subject center position is offset from the scan center by an amount exceeding a predefined range (e.g., 50 cm). Thus, the method can detect whether the center position of a subject is significantly offset from a scan center.\nIn accordance with a second aspect, there is provided a subject center position detection method comprising the steps of: acquiring a projection profile by imaging a subject from a direction orthogonal to a plane containing a body axis of the subject by an X-ray CT apparatus for performing a scan by an X-ray tube rotating around the subject; defining an X-ray detector channel index corresponding to a calculated centroid of the projection profile as a centroid channel index; examining a projection value in the projection profile toward a direction along a first direction of the subject from the centroid channel index; defining a channel index having a projection value no more than a predetermined projection value threshold Pth as a first threshold channel index; examining on the other hand a projection value in the projection profile toward a direction along a second direction of the subject opposite to the first direction from the centroid channel index; defining a channel index having a projection value no more than the predetermined projection value threshold Pth as a second threshold channel index; defining an intermediate value between the first and second threshold channel indices as a subject position channel index; and obtaining a subject center position from the subject position channel index.\nA projection profile acquired by imaging a subject lying on a top plate from a direction orthogonal to a plane containing the subject\"\"s body axis contains a projection of the top plate as well as a projection of the subject. Thus, the centroid of the projection profile is shifted toward the top plate relative to a subject center position. Therefore, the subject center position should be determined on a side opposite to the top plate with respect to the centroid of the projection profile.\nThe subject center position detection method of the second aspect introduces xe2x80x9cprojection value threshold Pthxe2x80x9d as a parameter, searches channels having respective values no more than the projection value threshold Pth in two directions, calculates an intermediate of the retrieved two channels, and determines the subject center position on a side opposite to the top plate with respect to the centroid of the projection profile based on the intermediate. Thus, the subject center position can be accurately determined with the effect of the top plate eliminated.\nIn the subject center position detection method of the above configuration, the projection value threshold Pth may be selected empirically as, for example, xc2xd of a projection value corresponding to the centroid channel index.\nMoreover, in the subject center position detection method of the above configuration, the subject center position can be determined more accurately by first removing profiles other than the subject\"\"s profile from the projection profile, and then calculating the centroid.\nIn accordance with a fifth aspect, there is provided an X-ray CT apparatus for performing a scan by an X-ray tube rotating around a subject, comprising: projection profile acquisition means for acquiring a projection profile by imaging the subject; subject position detection means for detecting a subject center position based on the acquired projection profile; and warning displaying means for displaying a warning when the distance between the subject center position and a scan center exceeds a predefined range.\nThe X-ray CT apparatus of the fifth aspect can suitably implement the X-ray CT imaging method as described regarding the first aspect.\nIn accordance with a sixth aspect, there is provided an X-ray CT apparatus for performing a scan by an X-ray tube rotating around a subject, comprising subject center position detection means configured to: define an X-ray detector channel index corresponding to a calculated centroid of a projection profile acquired by imaging a subject from a direction orthogonal to a plane containing a body axis of the subject as a centroid channel index; examine a projection value in the projection profile toward a direction along a first direction of the subject from the centroid channel index; define a channel index having a projection value no more than a predetermined projection value threshold Pth as a first threshold channel index; examine on the other hand a projection value in the projection profile toward a direction along a second direction of the subject opposite to the first direction from the centroid channel index; define a channel index having a projection value no more than the predetermined projection value threshold Pth as a second threshold channel index; define an intermediate value between the first and second threshold channel indices as a subject position channel index; and obtain a subject center position from the subject position channel index.\nThe X-ray CT apparatus of the sixth aspect can suitably implement the subject center position detection method as described regarding the second aspect.\nIn the X-ray CT apparatus of the above configuration, the projection value threshold Pth may be selected empirically as, for example, xc2xd of a projection value corresponding to the centroid channel index.\nMoreover, in the X-ray CT apparatus of the above configuration, the subject center position can be determined more accurately by first removing profiles other than the subject\"\"s profile from the projection profile, and then calculating the centroid.\nTo remove profiles other than the subject\"\"s profile from the projection profile, a one-dimensional minimum value filter, followed by a one-dimensional maximum value filter, for example, may be applied to the projection file.\nThe X-ray CT imaging method, subject center position detection method and X-ray CT apparatus of the present invention can detect whether the center position of a subject is offset from a scan center, and can notify an operator of the result of the detection. Thus, even when the subject center position is significantly offset from the scan center, the position of the subject is adjusted and imaging is repeated, whereby an X-ray image with good quality can be obtained.\nFurther objects and advantages of the present invention will be apparent from the following description of the preferred embodiments of the invention as illustrated in the accompanying drawings."}
{"text": "Spray dryers basically have three essential functions of (1) spraying feed liquid, (2) drying resulting microdrops and (3) separating and recovering a resulting fine powder product, and ordinarily have an atomizer, a drying chamber and a fine powder recovering unit, respectively, corresponding to the above functions.\nAs examples of these spray dryers, there have hitherto been known those shown in FIGS. 9 and 10 (see Japanese Patent Publication No. 32601/1983 and Japanese Utility Model Laid-Open No. 26950/1983).\nIn the spray dryer of FIG. 9, a feed liquid is sprayed into a drying chamber 1 from a nozzle 5 and is momentarily heated by hot air from an inlet 2, whereby the liquid component of the feed liquid is evaporated and the solid component is made into a fine powder. The majority of the fine powder is taken out as a product via a rotary valve 6. Part of the fine powder accompanies hot air, passes a waste gas outlet 3 and enters a cyclone 7, and is recovered from a rotary valve 8.\nThe spray dryer of FIG. 10 is a type using, as an atomizer, a rotary spraying disc (a rotary disc) 10 in place of the nozzle 5.\nIn the spray dryers as shown in FIGS. 9 and 10, a fine powder product is recovered separately in the main section of the dryer and in the cyclone. In this case, a powder of lighter particles having smaller particle diameters is recovered in the cyclone and a powder of heavier particles having larger particle diameters is taken out from the main section; accordingly, slight compositional variations in the product occurs.\nIn the past, this compositional variation posed little problem depending upon the type of product.\nIt has been required recently that spray dryers can perform, instead of conventional few grade, large quantity type production, many grade, small quantity type production especially for medicines, fine ceramics, etc.\nAlso, it is increasingly required that spray dryers produce a product of higher purity.\nUnder these new situations requiring a spray dryer capable of providing many grade, small quantity type production and a powder of the higher purity, the abovementioned compositional variation of product which posed little problem in the past has entered a new phase. Moreover a new problem has arisen.\nThat is, with the increase in the frequency of product switch-over, the rapid and complete washing of the deposit adhering to the inner wall of the drying chamber has become necessary. It is necessary because if washing takes a long time it is detrimental to the efficient execution of many grade, small quantity type production and also because the complete washing of the deposit adhering to the inner wall of the drying chamber, which is required in the product switchover when a product of higher purity and higher quality is needed, is very difficult. Moreover, when a product harmful to humans is produced, washing is extremely difficult because the operator is prohibited from making direct contact with the product.\nFurther, the requirement for a product of higher purity has turned the above-mentioned slight compositional variation among products into a problem."}
{"text": "1. Field of the Invention\nThe present invention relates to a computer terminal apparatus used for voting at an election, and an electronic system for voting at an election and totaling the votes cast.\n2. Description of the Related Art\nConventionally, at presidential elections, elections of Diet members, and the like, voters go to polling stations, where the voters write candidate names on ballot papers and cast the ballot papers into ballot boxes. Thereafter, vote counting and totaling of election returns are manually performed.\nAs described above, in the conventional voting at elections, voting operations, totaling operations, and the like are not automated much. Therefore, it takes much time to perform voting and totaling operations. In addition, vote totaling errors often occur."}
{"text": "1. Field of the Invention\nThis invention relates in general to an article of furniture, and, more particularly, to that class known as wardrobes.\n2. Description of the Related Art\nWardrobes must act as repositories for a large number of small and large wardrobe articles such as shirts, belts and jewelry. The problem is then to store these variously sized items in a compact, yet readily accessible manner for presentation and use that does not wrinkle or mess the items for wear."}
{"text": "1. Field of Invention\nThe present invention is directed to a novel bar-code symbol driven Internet Access Terminal, and method of accessing Internet-based information resources by scanning bar code symbols encoded with Uniform Resource Locators (URLs) or Domain Name/Path Name (DN/PN) portions thereof indicating the location of such information resources on the Internet.\n2. Background of the Invention\nPresently, several techniques have been developed for connecting to or accessing (i.e., xe2x80x9csurfingxe2x80x9d among) Internet-based information resources on the World Wide Web (WWW) using a conventional graphical user interface (GUI) based Internet browser programs, such as the Navigator(copyright) from Netscape Communications, Inc. or the Internet Explorer(copyright) from Microsoft, Inc. The availability of any particular Web-site surfing technique depends on where the Internet user finds him or her self in the Internet browser program at any particular instant of time.\nFor example, if one is currently at a Web-site (i.e., in a particular HyperText Markup Language (HTML) document), at which there is a highlighted or embedded xe2x80x9clinkxe2x80x9d specifying the Internet address of another Web-site (i.e., another HTML document), then the user can access (i.e., xe2x80x9csurf toxe2x80x9d) this other Web-site by simply xe2x80x9cclicking onxe2x80x9d or selecting the highlighted URL with his or her pointing device (i.e., xe2x80x9cmousexe2x80x9d) in a conventional manner. The ability to connect to other Web-sites by simply pointing and clicking on a highlighted URLs embedded in an HTML document has contributed enormously to the growth and popularity of the Internet in the last few years.\nIf the user desires to connect to a Web-site that is not referenced in a HTML document viewed through a browser program, then the user oftentimes finds it necessary to manually enter the URL of the Web-site (e.g., xe2x80x9chttp://www.metrologic.comxe2x80x9d) into the browser program. This requires manually pressing a sequence of keys on a keyboard or remote control device, corresponding to the characters of the URL being selected. In addition to being time consuming, this Web-site surfing technique is prone to errors, causing the browser program to connect to the wrong Web-site or return a message stating that the Domain Name of the Web-site sought after cannot be found. From a practical point of view, this is quite frustrating to the Internet user.\nMore recently, with the growth and ever increasing complexity of the WWW, it is becoming popular to use printed publications (e.g., magazines, catalogues, directories, etc.) which list Web-sites, corresponding URLs and content descriptions, as navigational aids, much in same way that the TV GUIDE(copyright) magazine is used to help viewers select programs during television viewing. Again, however, each time a listed Web-site is to be viewed, the user is forced to manually enter into the Internet browser program, the lengthy character string associated with the URL of the Web-site being selected. Only thereafter, does the browser program automatically connect the user\"\"s client system (i.e., Internet Terminal) to the Internet Server computer supporting the selected Web-site.\nWhile the development of HTML and GUI-based Web browser programs have made accessing Web-sites relatively simple, the above shortcomings and drawbacks of prior art Internet technology clearly indicate a great need in the art for an improved system and method of accessing information resources on the Internet.\nAccordingly, it is a primary object of the present invention to provide an improved method and apparatus for accessing and navigating among information resources on the Internet, while avoiding the shortcomings and drawbacks of prior art systems and methodologies.\nA further object of the present invention is to provide an Internet Access System which includes a bar code symbol reader for reading DN/PN-encoded (Domain Name and Path Name encoded) or URL-encoded (Uniform Resource Location encoded) bar code symbols printed on various types of objects, including print media, which, when read thereby, automatically connects the Internet Access System to the Internet Server that contains the information resource specified by the scanned DN/PN-encoded or URL-encoded bar code symbol.\nA further object of the present invention is to provide such Internet Access System, wherein the bar code symbol reader may be a laser scanning bar code symbol reader, a CCD-type bar code symbol, or a Wand-type bar code symbol reader.\nA further object of the present invention is to provide such an Internet Access System, wherein the DN-encoded bar code symbol printed on various types of print media is a DN/PN-encoded truncated-type bar code symbol, having a very low height to length ratio, thereby allowing many URL-encoded bar code symbols to be printed on a single sheet or page of a Web-site guide, along with their corresponding human-readable URLs and content descriptions.\nA further object of the present invention is to provide such an Internet Access System, wherein the URL-encoded bar code symbol printed on various types of print media is a URL-encoded truncated-type bar code symbol, having a very low height to length ratio, thereby allowing many URL-encoded bar code symbols to be printed on a single sheet or page of a Web-site guide, along with their corresponding human-readable URLs and content descriptions.\nAnother object of the present invention is to provide such an Internet Access System in the form of a desktop, laptop or palmtop computer system that is connected to the Internet by way of an Internet Service Provider (ISP), wherein the computer system has a GUI-based web browser program and a programmed bar code symbol scanner interfaced therewith for automatically surfing to information resources (e.g., Web-sites) listed in a Web-site guide by simply scanning corresponding URL-encoded bar code symbols printed on the pages thereof.\nA further object of the present invention is to provide an Internet Access System in the form of an interactive web-based television system, wherein the web-based television system comprises an Internet terminal unit connected to the Internet by way of an ISP, an audio-visual display monitor for displaying graphical and audio information content of Web-sites, and a portable Internet surfing device having a wireless IR-based communication link to the Internet terminal unit and an integrated bar code symbol scanner for automatically surfing to (or among) Web-sites listed in a Web-site guide by simply scanning corresponding URL-encoded bar code symbols printed on the pages thereof.\nA further object of the present invention is to provide such an Internet Access System in the form of a Scanner Integrated Terminal, wherein the Scanner Integrated Terminal is connectable to the Internet by way of a wireless (RF) link to an ISP, and has an integrated GUI-based web browser program, display panel, keypad, and programmed bar code symbol scanner for automatically surfing to Web-sites listed in a Web-site guide by simply scanning corresponding DN-encoded (or URL-encoded) bar code symbols printed on the pages thereof.\nAnother object of the present invention is to provide a novel method of surfing to Web-sites on the Internet by scanning URL-encoded bar code symbols into GUI-based web browser programs, without the need of manual data entry operations or the like.\nA further object of the present invention is to provide an Internet Access System, which includes an optical character reader programmed to read the character strings of URLs printed on various types of print media which, when read thereby, automatically connects the Internet Client System to the Internet Server that supports the Web-site specified by the read URL.\nAnother object of the present invention is to provide such an Internet Access System in the form of a desktop, laptop or palmtop computer system that is connected to the Internet by way of an Internet Service Provider (ISP), wherein the computer system has a GUI-based web browser program and a programmed optical character reader interfaced therewith for automatically surfing to Web-sites listed in a Web-site guide by simply reading corresponding URLs printed on the pages thereof.\nA further object of the present invention is to provide an Internet Access System in the form of an interactive web-based television system, wherein the web-based television system comprises an Internet terminal unit connected to the Internet by way of an ISP, an audio-visual (AV) display monitor for displaying graphical and audio information content of Web-sites, and a portable Internet surfing device having a wireless IR-based communication link to the Internet Terminal unit and an integrated optical character reader for automatically surfing to Web-sites listed in a Web-site guide by simply scanning corresponding URLs printed on the pages thereof.\nA further object of the present invention is to provide such an Internet Access System in the form of a Scanner Integrated Terminal, wherein the Scanner Integrated Terminal is connectable to the Internet by way of a wireless (RF) link to an ISP, and has an integrated GUI-based web browser program, display panel, keypad, and programmed optical character reader for automatically surfing to Web-sites listed in a Web-site guide by simply scanning corresponding URLs printed on the pages thereof.\nAnother object of the present invention is to provide a novel method of surfing to Web-sites on the Internet by optically scanning the character strings of URLs into GUI-based web browser programs, without the need of manual data entry operations and the like.\nAnother object of the present invention is to provide a novel system and method for printing URL-encoded bar code symbols on various types of print media which, when read thereby, automatically connects the Internet Client System to the Internet Server that supports the Web-site specified by the scanned URL-encoded bar code symbol.\nAnother object of the present invention is to provide such a system in the form of suitably programmed desktop Web-linked computer workstation or laptop computer system having a graphical user interface and an Internet browser program having a plug-in type URL-Menu Composition/Printing Module that supports the bar code menu composition, transmission and printing processes of the present invention.\nA further object of the present invention is to provide such a system, in which the Internet browser program and plug-in URL-Menu Composition/Printing Module allow the user to easily compile the following information structure during operation of the Internet browser program: a list of human-readable URLs corresponding to a set or course of information serving Web-sites, a set of URL-encoded bar code symbol data structures corresponding thereto, a set of titles assigned to the URLs, a set of brief descriptions of content served at the Web-sites, and date or dates the set of Web-sites were last visited by the compiler for subsequent printing in the form of a menu.\nA further object of the present invention is to provide such a system, in which the compiled information structure can be electronically transmitted to a remote site, by e-mail, facsimile transmission, or other protocol available over the Internet, and thereafter printed out using appropriate print-driver software, so that the Web-sites listed in the compiled information structure can be easily visited by reading the corresponding URL-encoded bar code symbol into the Internet browser program using a bar code scanner.\nAnother object of the present invention is to provide a novel Web-site guide or directory comprising one or more sheets having a plurality of URL-encoded bar code symbols, corresponding human-readable URLs and Web-site content descriptions printed thereon.\nAnother object of the present invention is to provide a novel Web-based document tracking and management system, in which each printed document (e.g., brochure, drawing, 3-D objects or specimens, photograph, specification, blue-print, and the like) to be tracked and managed within the system is uniquely assigned and affixed with a printed URL-encoded bar code symbol that specifies (i.e.,points to) a particular information storage location within an HTML-encoded relational database management system (RDBMS), realized on an Web-based information server located on the Internet or behind the corporate firewall of an HTTP-supporting Intranet.\nAnother object of the present invention is to provide such a Web-based document tracking and management system, in which information relating to a document registered within the system can be easily accessed therefrom by simply scanning the printed URL-encoded bar code symbol on the document using the bar code symbol reader associated with an Internet Access Terminal (i.e., client computer system) of the present invention.\nAnother object of the present invention is to provide such a Web-based document tracking and management system, in which the Internet Access Terminal of each user can be used to: (1) assign a URL-encoded bar code symbol to any document to be registered with the system; (2) enter document-related information into the information storage location assigned to the document within the Web-based RDBMS; (3) print out URL-encoded bar code symbols for attachment to corresponding documents registered within the system; (4) modify at any time the document-related information currently stored in the information storage location assigned to the document within the Web-based RDBMS; and (5) access such document related information from the system by simply reading its corresponding URL-encoded bar code symbol.\nAnother object of the present invention to provide a Web-based package routing, tracking and delivering system and method that uses URL-encoded bar code symbols on parcels and packages.\nA further object of the present invention is to provide such a Web-based package routing, tracking and delivery system and method, wherein one or more Routing, Tracking and Delivery (RTD) Internet Server Subsystems are connected to the Internet and updated at any instant of time with package related information produced by either (i) a Package Log-In/Shipping Subsystem that is located at a product shipping location (e.g., warehouse) and connected to the RTD Internet Server by a first data communication link, (ii) a Package Routing Subsystem that is located at a hub station and connected to the RTD Internet Server by a second data communication link, or (iii) a Portable Package Delivery Subsystem that is carried by the package delivery person and connected to the RTD Internet Server by a wireless data communication link.\nA further object of the present invention is to provide such a Web-based package routing, tracking and delivery system and method, wherein at each remote hub station within the system: the URL/ZIP-CODE encoded bar code symbol is automatically scanned; the encoded destination Zip Code is locally recovered and used to route the package at the hub station; and the locally recovered URL is used to access the RTD Internet Server and update the location of the package within the system.\nAnother object of the present invention is to provide such a Web-based package routing, tracking and delivery system and method, wherein when the URL/ZIP-CODE Encoded bar code symbol is read by the deliveryman or courier using a portable (wireless) Internet Access Terminal, the recovered URL is used to automatically access the RTD Internet Server and display on the Internet Access Terminal, precise up-to-date parcel delivery information, payment-method terms, handling instructions and the like.\nThese and other Objects of the Present Invention will become apparent hereinafter."}
{"text": "The present invention relates to the field of macrolide compounds having antibacterial activity, pharmaceutical compositions containing the compounds, and methods of treating bacterial infections with the compounds.\nErythromycins are well-known antibacterial agents widely used to treat and prevent bacterial infection caused by Gram-positive and Gram-negative bacteria. However, due to their low stability in acidic environment, they often carry side effects such as poor and erratic oral absorption. As with other antibacterial agents, bacterial strains having resistance or insufficient susceptibility to erythromycin have developed over time and are identified in patients suffering from such ailments as community-acquired pneumonia, upper and lower respiratory tract infections, skin and soft tissue infections, meningitis, hospital-acquired lung infections, and bone and joint infections. Particularly problematic pathogens include methicillin-resistant Staphylococcus aureus (MRSA), vancomycin-resistant enterococci (VRE) and penicillin- and macrolide-resistant Streptococcus pneumoniae. Therefore, continuing efforts are called for to identify new erythromycin derivative compounds with improved antibacterial activity, and/or unanticipated selectivity against various target microorganisms, particularly erythromycin-resistant strains.\nThe following references relate to various erythromycin derivatives disclosed as having antibacterial activity:\nEP 216,169 and U.S. Pat. No. 4,826,820 to Brain et al. disclose antibacterially active 6-carbamate erythromycin derivatives stated to xe2x80x9chave antibacterial properties, in particular against Gram-positive bacteria but also against some Gram-negative bacteria.xe2x80x9d\nU.S. Pat. No. 5,444,051, No. 5,561,118, and No. 5,770,579, all to Agouridas et al., disclose erythromycin compounds such as those of the formulae \nwherein substituents are as described in the respective references, which are all stated to be useful as antibiotics.\nU.S. Pat. No. 5,866,549 to Or et al. and WO 98/09978 (Or et al.) disclose 6-O-substituted ketolides stated to have increased acid stability relative to erythromycin A and 6-O-methyl erythromycin A and enhanced activity toward gram negative bacteria and macrolide resistant gram positive bacteria.\nWO 97/17356 (Or et al.) discloses tricyclic erythromycin derivatives stated to be useful in the treatment and prevention of bacterial infections.\nWO 99/21871 (Phan et al.) discloses 2-halo-6-O-substituted ketolide derivatives of the formula \nwherein substituents are as described in the respective reference, which are stated to possess antibacterial activity.\nWO 99/21864 (Or et al.) discloses 6,11-bridged erythromycin derivatives having antibacterial activity.\nEP1146051 to Kaneko et al. discloses macrolide compounds of the following formula that are useful as antibacterial and antiprotozoal agents in mammals, \nwherein substituents are as described in the reference.\nWO 00/75156 (Phan et al.) discloses 6-O-carbamate ketolide derivatives that are useful as antibacterials for the treatment and prevention of infections in a mammal.\nThe invention provides compounds of Formula 1: \nWherein\nR1 and R2 are independently selected from hydrogen, optionally substituted xe2x80x94CRfRg(C1-C8)alkyl, optionally substituted xe2x80x94CRfRg(C2-C8)alkenyl, optionally substituted xe2x80x94CRfRg(C2-C8)alkynyl, optionally substituted cycloalkyl, and optionally substituted (C5-C7)-cycloalkenyl, provided that R1 and R2 are not both hydrogen, wherein the substituents are selected from halogen, alkyl, alkenyl, alkynyl, cycloalkyl, oxo, aryl, heteroaryl, heterocyclo, CN, nitro, xe2x80x94COORa, xe2x80x94OCORa, xe2x80x94ORa, xe2x80x94SRa, xe2x88x92SORa, xe2x80x94SO2Ra, xe2x80x94NRaRb, xe2x80x94CONRaRb, xe2x80x94OCONRaRb, xe2x80x94NHCORa, xe2x80x94NHCOORa, and xe2x80x94NHCONRaRb, wherein\nRa and Rb are independently selected from hydrogen, alkyl, alkenyl, alkynyl, cycloalkyl, aryl, heteroaryl, heterocyclo, aralkyl, heteroaralkyl, and heterocycloalkyl; and\nRf and Rg are independently selected from the group consisting of hydrogen, alkyl, alkenyl, alkynyl, cycloalkyl, aryl, heteroaryl, heterocyclo, COORa, and CONRaRb;\nor R1 and R2, together with the nitrogen atom to which they are attached, form an optionally substituted nitrogen-containing heterocycle, wherein the substituents are selected from halogen, alkyl, alkenyl, alkynyl, cycloalkyl, oxo, aryl, heteroaryl, heterocyclo, CN, nitro, xe2x80x94COORa, xe2x80x94OCORa, xe2x80x94ORa, xe2x80x94SRa, xe2x80x94SORa, xe2x80x94SO2Ra, xe2x80x94NRaRb, xe2x80x94CONRaRb, xe2x80x94OCONRaRb, xe2x80x94NHCORa, xe2x80x94NHCOORa, and xe2x80x94NHCONRaRb;\nR3 is hydrogen or xe2x80x94Wxe2x80x94V, wherein\nW is selected from the group consisting of (a) xe2x80x94NHxe2x80x94(CH2)pxe2x80x94, (b) xe2x80x94(CH2)qxe2x80x94, (c) xe2x80x94Oxe2x80x94(CH2)rxe2x80x94, (d) xe2x80x94NHxe2x80x94C1-C6alkenyl-, (e) xe2x80x94C1-C6alkenyl-, (f) xe2x80x94Oxe2x80x94C1-C6alkenyl-, (g) xe2x80x94NHxe2x80x94C1-C6alkynyl-, (h) xe2x80x94C1-C6alkynyl-, and (i) xe2x80x94Oxe2x80x94C1-C6alkynyl-;\np is 0 to 5;\nq is 0 to 5; and\nr is 0 to 5;\nand\nV is selected from the group consisting of (a) hydrogen, (b) aryl, (c) substituted aryl, (d) heteroaryl, (e) substituted heteroaryl, and (f) Ar1xe2x80x94Ar2, wherein Ar1 and Ar2 are independently selected from the group consisting of (i) aryl, (ii) substituted aryl, (iii) heteroaryl, and (iv) substituted heteroaryl;\nR4 is selected from hydrogen, halogen, and hydroxy;\nR5 is hydrogen or a hydroxy protecting group;\nR6 is selected from hydrogen, alkyl, C2-C10-alkenyl, C2-C10-alkynyl, aryl, heteroaryl, heterocyclo, aryl(C1-C10)alkyl, aryl(C2-C10)alkenyl, aryl(C2-C10)alkynyl, heterocyclo(C1-C10)alkyl, heterocyclo(C2-C10)alkenyl, and heterocyclo(C2-C10)alkynyl, C3-C6-cycloalkyl, C5-C8-cycloalkenyl, alkoxyalkyl containing 1-6 carbon atoms in each alkyl or alkoxy group, and alkylthioalkyl containing 1-6 carbon atoms in each alkyl or thioalkyl group;\nX and Xxe2x80x2, together with the carbon atom to which they are attached, form Cxe2x95x90O, Cxe2x95x90NRc, or Cxe2x95x90NORc, wherein Rc is independently selected from hydrogen, alkyl, alkenyl and alkynyl; and\nY and Yxe2x80x2, together with the carbon atom to which they are attached, form Cxe2x95x90O, xe2x80x94CHOH, Cxe2x95x90NRc, or Cxe2x95x90NORc, wherein Rc is independently selected from hydrogen, alkyl, alkenyl and alkynyl;\nor an optical isomer, enantiomer, diastereomer, racemate or racemic mixture thereof, or a pharmaceutically acceptable salt, esters or pro-drugs thereof.\nCompounds of the above formula are useful as antibacterial agents for the treatment of bacterial infections in a subject such as human and animal.\nThe present invention is also directed to a method of treating a subject having a condition caused by or contributed to by bacterial infection, which comprises administering to said subject a therapeutically effective amount of the compound of Formula 1.\nThe present invention is further directed to a method of preventing a subject from suffering from a condition caused by or contributed to by bacterial infection, which comprises administering to the subject a prophylactically effective amount of the compound of Formula 1.\nOther objects and advantages will become apparent to those skilled in the art from a review of the ensuing specification.\nRelative to the above description, certain definitions apply as follows.\nUnless otherwise noted, under standard nomenclature used throughout this disclosure the terminal portion of the designated side chain is described first, followed by the adjacent functionality toward the point of attachment.\nUnless specified otherwise, the terms xe2x80x9calkylxe2x80x9d, xe2x80x9calkenylxe2x80x9d, and xe2x80x9calkynyl,xe2x80x9d whether used alone or as part of a substituent group, include straight and branched chains having 1 to 8 carbon atoms, or any number within this range. The term xe2x80x9calkylxe2x80x9d refers to straight or branched chain hydrocarbons. xe2x80x9cAlkenylxe2x80x9d refers to a straight or branched chain hydrocarbon with at least one carbonxe2x80x94carbon double bond. xe2x80x9cAlkynylxe2x80x9d refers to a straight or branched chain hydrocarbon with at least one carbonxe2x80x94carbon triple bound. For example, alkyl radicals include methyl, ethyl, n-propyl, isopropyl, n-butyl, isobutyl, sec-butyl, t-butyl, n-pentyl, 3-(2-methyl)butyl, 2-pentyl, 2-methylbutyl, neopentyl, n-hexyl, 2-hexyl and 2-methylpentyl. xe2x80x9cAlkoxyxe2x80x9d radicals are oxygen ethers formed from the previously described straight or branched chain alkyl groups. xe2x80x9cCycloalkylxe2x80x9d groups contain 3 to 8 ring carbons and preferably 5 to 7 ring carbons. xe2x80x9cCycloalkenylxe2x80x9d groups contain 5 to 8 ring carbons and at least one carbonxe2x80x94carbon double bond. The alkyl, alkenyl, alkynyl, cycloalkyl, cycloalkenyl, and alkoxy group may be independently substituted with one or more members of the group including, but not limited to, halogen, alkyl, alkenyl, alkynyl, cycloalkyl, alkoxy, oxo, aryl, heteroaryl, heterocyclo, CN, nitro, xe2x80x94OCORa, xe2x80x94ORa, xe2x80x94SRa, xe2x80x94SORa, xe2x80x94SO2Ra, xe2x80x94COORa, xe2x80x94NRaRb, xe2x80x94CONaRb, xe2x80x94OCONRaRb, xe2x80x94NHCORa, xe2x80x94NHCOORa, and xe2x80x94NHCONRaRb, wherein Ra and Rb are independently selected from H, alkyl, alkenyl, alkynyl, cycloalkyl, aryl, heteroaryl, heterocyclo, aralkyl, heteroaralkyl, and heterocycloalkyl. xe2x80x9cAralkyl,xe2x80x9d xe2x80x9cheteroaralkyl,xe2x80x9d and xe2x80x9cheterocycloalkylxe2x80x9d are alkyl groups substituted with aryl, heteroaryl, and heterocyclo, respectively. xe2x80x9cArylalkenyl,xe2x80x9d xe2x80x9cheteroarylalkenyl,xe2x80x9d and xe2x80x9cheterocycloalkenylxe2x80x9d are alkenyl groups substituted with aryl, heteroaryl, and heterocyclo, respectively. xe2x80x9cArylalkynyl,xe2x80x9d xe2x80x9cheteroarylalkynyl,xe2x80x9d and xe2x80x9cheterocycloalkynylxe2x80x9d are alkynyl groups substituted with aryl, heteroaryl, and heterocyclo, respectively.\nThe term xe2x80x9cacylxe2x80x9d as used herein, whether used alone or as part of a substituent group, means an organic radical having 2 to 6 carbon atoms (branched or straight chain) derived from an organic acid by removal of the hydroxyl group. The term xe2x80x9cAcxe2x80x9d as used herein, whether used alone or as part of a substituent group, means acetyl.\nThe term xe2x80x9chaloxe2x80x9d or xe2x80x9chalogenxe2x80x9d means fluoro, chloro, bromo and iodo. (Mono-, di-, tri-, and per-)halo-alkyl is an alkyl radical substituted by independent replacement of the hydrogen atoms thereon with halogen.\nxe2x80x9cArylxe2x80x9d or xe2x80x9cAr,xe2x80x9d whether used alone or as part of a substituent group, is a carbocyclic aromatic radical including, but not limited to, phenyl, 1- or 2-naphthyl and the like. The carbocyclic aromatic radical may be substituted by independent replacement of 1 to 3 of the hydrogen atoms thereon with halogen, OH, CN, mercapto, nitro, amino, C1-C8-alkyl, aryl, heteroaryl, heterocyclo, C1-C8-alkoxyl, C1-C8-alkylthio, C1-C8-alkyl-amino, di(C1-C8-alkyl)amino, (mono-, di-, tri-, and per-) halo-alkyl, formyl, carboxy, alkoxycarbonyl, C1-C8-alkyl-COxe2x80x94Oxe2x80x94, C1-C8-alkyl-COxe2x80x94NHxe2x80x94, or carboxamide. Illustrative aryl radicals include, for example, phenyl, naphthyl, biphenyl, fluorophenyl, difluorophenyl, benzyl, benzoyloxyphenyl, carboethoxyphenyl, acetylphenyl, ethoxyphenyl, phenoxyphenyl, hydroxyphenyl, carboxyphenyl, trifluoromethylphenyl, methoxyethylphenyl, acetamidophenyl, tolyl, xylyl, dimethylcarbamylphenyl and the like. xe2x80x9cPhxe2x80x9d or xe2x80x9cPHxe2x80x9d denotes phenyl.\nWhether used alone or as part of a substituent group, xe2x80x9cheteroarylxe2x80x9d refers to a cyclic, fully unsaturated radical having from five to ten ring atoms of which one ring atom is selected from S, O, and N; 0-3 ring atoms are additional heteroatoms independently selected from S, O, and N; and the remaining ring atoms are carbon. The radical may be joined to the rest of the molecule via any of the ring atoms. Exemplary heteroaryl groups include, for example, pyridinyl, pyrazinyl, pyrimidinyl, pyridazinyl, pyrroyl, pyrazolyl, imidazolyl, thiazolyl, oxazolyl, isoxazolyl, thiadiazolyl, triazolyl, triazinyl, oxadiazolyl, thienyl, furanyl, quinolinyl, isoquinolinyl, indolyl, isothiazolyl, N-oxo-pyridyl, 1,1-dioxothienyl, benzothiazolyl, benzoxazolyl, benzothienyl, quinolinyl-N-oxide, benzimidazolyl, benzopyranyl, benzisothiazolyl, benzisoxazolyl, benzodiazinyl, benzofurazanyl, indazolyl, indolizinyl, benzofuryl, cinnolinyl, quinoxalinyl, pyrrolopyridinyl, furopyridinyl (such as furo[2,3-c]pyridinyl, furo[3,2-b]pyridinyl, or furo[2,3-b]pyridinyl), imidazopyridinyl (such as imidazo[4,5-b]pyridinyl or imidazo[4,5-c]pyridinyl), naphthyridinyl, phthalazinyl, purinyl, pyridopyridyl, quinazolinyl, thienofuryl, thienopyridyl, and thienothienyl. The heteroaryl group may be substituted by independent replacement of 1 to 3 of the hydrogen atoms thereon with halogen, OH, CN, mercapto, nitro, amino, C1-C8-alkyl, aryl, heteroaryl, heterocyclo, C1-C8-alkoxyl, C1-C8-alkylthio, C1-C8-alkyl-amino, di(C1-C8-alkyl)amino, (mono-, di-, tri-, and per-) halo-alkyl, formyl, carboxy, alkoxycarbonyl, C1-C8-alkyl-COxe2x80x94Oxe2x80x94, C1-C8-alkyl-COxe2x80x94NHxe2x80x94, or carboxamide. Heteroaryl may be substituted with a mono-oxo to give for example a 4-oxo-1 H-quinoline.\nThe terms xe2x80x9cheterocycle,xe2x80x9d xe2x80x9cheterocyclic,xe2x80x9d and xe2x80x9cheterocycloxe2x80x9d refer to an optionally substituted, fully saturated, partially saturated, or non-aromatic cyclic group which is, for example, a 3- to 7-membered monocyclic, 7- to 11-membered bicyclic, or 10- to 15-membered tricyclic ring system, which has at (east one heteroatom in at least one carbon atom containing ring. Each ring of the heterocyclic group containing a heteroatom may have 1, 2, or 3 heteroatoms selected from nitrogen atoms, oxygen atoms, and sulfur atoms, where the nitrogen and sulfur heteroatoms may also optionally be oxidized. The nitrogen atoms may optionally be quaternized. The heterocyclic group may be attached at any heteroatom or carbon atom.\nExemplary monocyclic heterocyclic groups include pyrrolidinyl; oxetanyl; pyrazolinyl; imidazolinyl; imidazolidinyl; oxazolinyl; oxazolidinyl; isoxazolinyl; thiazolidinyl; isothiazolidinyl; tetrahydrofuryl; piperidinyl; piperazinyl; 2-oxopiperazinyl; 2-oxopiperidinyl; 2-oxopyrrolidinyl; 4-piperidonyl; tetrahydropyranyl; tetrahydrothiopyranyl; tetrahydrothiopyranyl sulfone; morpholinyl; thiomorpholinyl; thiomorpholinyl sulfoxide; thiomorpholinyl sulfone; 1,3-dioxolane; dioxanyl; thietanyl; thiiranyl; 2-oxazepinyl; azepinyl; and the like. Exemplary bicyclic heterocyclic groups include quinuclidinyl; tetrahydroisoquinolinyl; dihydroisoindolyl; dihydroquinazolinyl (such as 3,4-dihydro-4-oxo-quinazolinyl); dihydrobenzofuryl; dihydrobenzothienyl; benzothiopyranyl; dihydrobenzothiopyranyl; dihydrobenzothiopyranyl sulfone; benzopyranyl; dihydrobenzopyranyl; indolinyl; chromonyl; coumarinyl; isochromanyl; isoindolinyl; piperonyl; tetrahydroquinolinyl; and the like. The heterocyclic group may be substituted by independent replacement of 1 to 3 hydrogen atoms thereon with OH, CN, mercapto, nitro, amino, C1-C8-alkyl, aryl, heteroaryl, heterocyclo, C1-C8-alkoxyl, C1-C8-alkylthio, C,-C8-alkyl-amino, di(C1-C8-alkyl)amino, (mono-, di-, tri-, and per-) halo-alkyl, formyl, carboxy, alkoxycarbonyl, C1-C8-alkyl-COxe2x80x94Oxe2x80x94, C1-C8-alkyl-COxe2x80x94NHxe2x80x94, or carboxamide.\nSubstituted aryl, substituted heteroaryl, and substituted heterocycle may also be substituted with a second substituted-aryl, a second substituted-heteroaryl, or a second substituted-heterocycle to give, for example, a 4-pyrazol-1-yl-phenyl or 4-pyridin-2-yl-phenyl.\nDesignated numbers of carbon atoms (e.g., C1-8) shall refer independently to the number of carbon atoms in an alkyl or cycloalkyl moiety or to the alkyl portion of a larger substituent in which alkyl appears as its prefix root.\nUnless specified otherwise, it is intended that the definition of any substituent or variable at a particular location in a molecule be independent of its definitions elsewhere in that molecule. It is understood that substituents and substitution patterns on the compounds of this invention can be selected by one of ordinary skill in the art to provide compounds that are chemically stable and that can be readily synthesized by techniques known in the art as well as those methods set forth herein.\nThe term xe2x80x9chydroxy protecting groupxe2x80x9d refers to groups known in the art for such purpose. Commonly used hydroxy protecting groups are disclosed, for example, in T. H. Greene and P. G. M. Wuts, Protective Groups in Organic Synthesis, 2nd edition, John Wiley and Sons, New York (1991), which is incorporated herein by reference. Illustrative hydroxyl protecting groups include but are not limited to tetrahydropyranyl; benzyl; methylthiomethyl; ethythiomethyl; phenylsulfonyl; triphenylmethyl; trisubstituted silyl such as trimethyl silyl, triethylsilyl, tributylsilyl, tri-isopropylsilyl, t-butyldimethylsilyl, tri-t-butylsilyl, methyidiphenylsilyl, ethyidiphenylsilyl, and t-butyidiphenylsilyl; acyl and aroyl such as acetyl, pivaloylbenzoyl, 4-methoxybenzoyl, and 4-nitrobenzoyl; and alkoxycarbonyl such as methoxycarbonyl, ethoxycarbonyl, benzyloxycarbonyl.\nWhere the compounds according to this invention have at least one stereogenic center, they may accordingly exist as enantiomers. Where the compounds possess two or more stereogenic centers, they may additionally exist as diastereomers. Furthermore, some of the crystalline forms for the compounds may exist as polymorphs and as such are intended to be included in the present invention. In addition, some of the compounds may form solvates with water (i.e., hydrates) or common organic solvents, and such solvates are also intended to be encompassed within the scope of this invention.\nSome of the compounds of the present invention may have trans and cis isomers. In addition, where the processes for the preparation of the compounds according to the invention give rise to mixture of stereoisomers, these isomers may be separated by conventional techniques such as preparative chromatography. The compounds may be prepared as a single stereoisomer or in racemic form as a mixture of some possible stereoisomers. The non-racemic forms may be obtained by either synthesis or resolution. The compounds may, for example, be resolved into their component enantiomers by standard techniques, such as the formation of diastereomeric pairs by salt formation. The compounds may also be resolved by covalent linkage to a chiral auxiliary, followed by chromatographic separation and/or crystallographic separation, and removal of the chiral auxiliary. Alternatively, the compounds may be resolved using chiral chromatography.\nThe phrase xe2x80x9ca pharmaceutically acceptable saltxe2x80x9d denotes one or more salts of the free base which possess the desired pharmacological activity of the free base and which are neither biologically nor otherwise undesirable. These salts may be derived from inorganic or organic acids. Examples of inorganic acids are hydrochloric acid, nitric acid, hydrobromic acid, sulfuric acid, or phosphoric acid. Examples of organic acids are acetic acid, propionic acid, glycolic acid, lactic acid, pyruvic acid, malonic acid, succinic acid, malic acid, maleic acid, fumaric acid, tartaric acid, citric acid, benzoic acid, cinnamic acid, mandelic acid, methanesulfonic acid, ethanesulfonic acid, p-toluenesulfonic acid, salicyclic acid and the like. Suitable salts are furthermore those of inorganic or organic bases, such as KOH, NaOH, Ca(OH)2, Al(OH)3, piperidine, morpholine, ethylamine, triethylamine and the like.\nIncluded within the scope of the invention are the hydrated forms of the compounds which contain various amounts of water, for instance, the hydrate, hemihydrate, and sesquihydrate forms. The present invention also includes within its scope prodrugs of the compounds of this invention. In general, such prodrugs will be functional derivatives of the compounds which are readily convertible in vivo into the required compound. Thus, in the methods of treatment of the present invention, the term xe2x80x9cadministeringxe2x80x9d shall encompass the treatment of the various disorders described with the compound specifically disclosed or with a compound which may not be specifically disclosed, but which converts to the specified compound in vivo after administration to the patient. Conventional procedures for the selection and preparation of suitable prodrug derivatives are described, for example, in xe2x80x9cDesign of Prodrugsxe2x80x9d, ed. H. Bundgaard, Elsevier, 1985.\nThe term xe2x80x9csubjectxe2x80x9d includes, without limitation, any animal or artificially modified animal. As a particular embodiment, the subject is a human.\nThe term xe2x80x9cdrug-resistantxe2x80x9d or xe2x80x9cdrug-resistancexe2x80x9d refers to the characteristics of a microbe to survive in presence of a currently available antimicrobial agent such as an antibiotic at its routine, effective concentration.\nThe compounds described in the present invention possess antibacterial activity due to their novel structure, and are useful as antibacterial agents for the treatment of bacterial infections in humans and animals.\nIn particular, compounds of Formula 1 wherein R1 and R2 are independently selected from hydrogen, substituted C1-C8-alkyl, substituted xe2x80x94CH2C2-C8-alkenyl, and substituted xe2x80x94CH2C2-C8-alkynyl, wherein the substituents are selected from CN, nitro, xe2x80x94COORa, xe2x80x94OCORa, xe2x80x94ORa, xe2x80x94SRa, xe2x80x94SORa, xe2x80x94SO2Ra, xe2x80x94NRaRb, xe2x80x94CONaRb, xe2x80x94OCONRaRb, xe2x80x94NHCORa, xe2x80x94NHCOORa, and xe2x80x94NHCONRaRb, wherein Ra and Rb are independently selected from hydrogen, alkyl, alkenyl, alkynyl, cycloalkyl, aryl, heteroaryl, heterocyclo, aralkyl, heteroaralkyl, and heterocycloalkyl; provided that R1 and R2 are not both hydrogen;\nR3 is hydrogen;\nR4 is selected from hydrogen, halogen, and hydroxy;\nR5 is hydrogen or a hydroxy protecting group;\nR6 is selected from hydrogen, alkyl, C2-C10-alkenyl, C2-C10-alkynyl, aryl, heteroaryl, heterocyclo, aryl(C,-C10)alkyl, aryl(C2-C10)alkenyl, aryl(C2-C10)alkynyl, heterocyclo(C1-C10)alkyl, heterocyclo(C2-C10)alkenyl, and heterocyclo(C2-C10)alkynyl, C3-C6-cycloalkyl, C5-C8-cycloalkenyl, alkoxyalkyl containing 1-6 carbon atoms in each alkyl or alkoxy group, and alkylthioalkyl containing 1-6 carbon atoms in each alkyl or thioalkyl group;\nX and Xxe2x80x2, together with the carbon atom to which they are attached, form Cxe2x95x90O, Cxe2x95x90NRc, or Cxe2x95x90NORc, wherein Rc is independently selected from hydrogen, alkyl, alkenyl and alkynyl; and\nY and Yxe2x80x2, together with the carbon atom to which they are attached, form Cxe2x95x90O, xe2x80x94CHOH, Cxe2x95x90NRc, or Cxe2x95x90NORc, wherein Rc is independently selected from hydrogen, alkyl, alkenyl and alkynyl\nare embodiments of the present invention for such purposes. More particularly, R2 is hydrogen, R4 is hydrogen or fluorine, X and Xxe2x80x2 form Cxe2x95x90O together with the carbon atom to which they are attached, and Y and Yxe2x80x2 form Cxe2x95x90O together with the carbon atom to which they are attached.\nCompounds of Formula 1, which are represented by Formula 1xe2x80x2: \nwherein R1, R2, R3, and R4 are as described above are also embodiments of this invention. More particularly, R2 and R3 are hydrogen and R4 is fluorine.\nCompounds of Formula 1, which are represented by Formula 1xe2x80x3: \nwherein R1, R2, R3, and R4 are as described above, are further embodiments of this invention. More particularly, R2 and R3 are hydrogen and R4 is fluorine.\nCompounds of Formula 1 wherein R5 may be selected from acyl and aroyl are further embodiments of this invention.\nThis invention also provides processes for preparing the instant compounds.\nThe compounds of Formula I may be prepared from readily available starting materials such as erythromycin and erythromycin derivatives well known in the art. Outlined in Schemes 1 through 18 are representative procedures to prepare the compounds of the instant invention: \nScheme I illustrates the method of synthesis of the 2xe2x80x2,4xe2x80x3-diacetyl-6-carbamyl-11,12-dideoxy-11,12-iminocarbonyloxyerythromycin A (VI) and the 2xe2x80x2-acetyl-6-carbamyl-11,12-dideoxy-3-O-descladinosyl-11,12-iminocarbonyloxyerythromycin A (1a) precursors to the compounds of the invention.\nErythromycin A is treated with acetic anhydride in the presence of a tertiary amine base, such as triethylamine, diisopropylethylamine, or pyridine, and an acylation catalyst, such as DMAP, in a suitable solvent such as methylene chloride, chloroform or THF at a temperature ranging from xe2x88x9220xc2x0 C. to 37xc2x0 C. for 2 to 48 hours to afford 2xe2x80x2,4xe2x80x3,11-triacetylerythromycin A (I). The 10,11-anhydro derivative (II) can be readily obtained by treatment of I with a base in an inert solvent such as THF, dioxane, DME, or DMF at a temperature ranging from xe2x88x9278xc2x0 C. to 80xc2x0 C. for 1-24 hours. Suitable bases to effect the elimination reaction include, but are not limited to, sodium hexamethyidisilazide, potassium hexamethyldisilazide, LDA, lithium tetramethylpiperidide, DBU, and tetramethylguanidine. It will be apparent to one skilled in the art that alternative methods for synthesis of 2xe2x80x2,4xe2x80x3-diacetyl-10,11-anhydroerythromycin A are available, including conversion of erythromycin A to the 11,12-cyclic carbonate derivative with ethylene carbonate, followed by elimination with tetramethylguanidine, as described in Hauske, J. R. and Kostek, G., J. Org. Chem. 1982, 47,1595. Selective protection of the 2xe2x80x2 and 4xe2x80x3-hydroxyl groups can then be readily accomplished with acetic anhydride in the presence of a tertiary amine base. Likewise, alternative protecting group strategies may be employed. For example, erythromycin A may be treated with benzoic anhydride, propionic anhydride, or formic acetic anhydride under similar conditions as described above to obtain the 2xe2x80x2,4xe2x80x3,11-triacylated erythromycin A derivative followed by elimination to afford the corresponding 10,11-anhydro compound.\nOnce the suitably protected 10,11-anhydro derivative is obtained, derivatization of both tertiary hydroxyl groups can be carried out by treatment with trichloroacetylisocyanate in an inert solvent, such as methylene chloride, chloroform, or THF at a temperature ranging from xe2x88x9220xc2x0 C. to 37xc2x0 C. for 1-24 hours to yield the di-(N-trichloroacetyl)carbamate derivative (III). The N-trichloroacetylcarbamate functionalities can be hydrolyzed to the corresponding primary carbamates by treatment with a suitable base, such as triethylamine, in an aqueous solvent mixture, such as methanol/water for 1-24 hours at a temperature ranging from 20xc2x0 C. to 80xc2x0 C. Alternative bases may likewise be used to effect this conversion, such as sodium hydroxide, potassium hydroxide, sodium carbonate and potassium carbonate. Under the reaction conditions, the primary carbamate formed at the 12-position undergoes spontaneous Michael addition to the electrophilic 11-position of the xcex1,xcex2-unsaturated ketone and the 2xe2x80x2-acetoxy group is hydrolyzed to the corresponding hydroxyl to afford the cyclic carbamate derivative (IV). Compound IV is generally isolated as a mixture of methyl epimers at the C10-position, which can be readily converted to the desired C10-xcex2-methyl epimer (V) by treatment with an equilibrating base, such as potassium t-butoxide, tetramethylguanidine, or DBU in a suitable solvent, such as THF, dioxane, DME, DMF or t-butanol at a temperature ranging from xe2x88x9278xc2x0 C. to 80xc2x0 C. for 1 to 24 hours. Reprotection of the 2xe2x80x2-hydroxyl group to give VI can be carried out by treatment with acetic anhydride in the presence of a tertiary amine base, such as triethylamine, diisopropylethylamine, or pyridine, and optionally an acylation catalyst, such as DMAP, in a suitable solvent such as methylene chloride, chloroform or THF at a temperature ranging from xe2x88x9220xc2x0 C. to 37xc2x0 C. for 2 to 48 hours. It is understood that an orthogonal protection strategy of the sugar hydroxyls may also be employed by treatment of V with alternate reagents such as benzoic anhydride, benzyl chloroformate, hexamethyidisilazane, or a trialkylsilyl chloride. Finally, selective removal of the cladinose sugar can be accomplished by reaction of VI with an acid, such as hydrochloric, sulfuric, chloroacetic, and trifluoroacetic, in the presence of alcohol and water to afford 1a. Reaction time is typically 0.5-24 hours at a temperature ranging from xe2x88x9210xc2x0 C. to 37xc2x0 C. \nScheme 2 depicts synthesis of compounds of formulae 1b, 1c and 1d, wherein RCHO is an aldehyde (R may be a member of the group including, but not limited to, hydrogen, alkyl, alkenyl, alkynyl, cycloalkyl, aryl, heteroaryl, heterocycle, arylalkenyl, arylalkynyl, aralkyl, heteroarylalkenyl, heteroarylalkynyl, heteroarylalkyl, heterocycloalkenyl, heterocycloalkynyl, and heterocycloalkyl). Oxidation of the 3-hydroxy group of 1a to yield compound 1b can be effected with DMSO and a carbodiimide, such as EDCI, in the presence of pyridinium trifluoroacetate in a suitable solvent, such as methylene chloride, for 1 to 24 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 37xc2x0 C. Alternative methods of oxidation include N-chlorosuccinimide and dimethylsulfide complex followed by treatment with a tertiary amine base, Dess-Martin periodinane, or oxalyl chloride/DMSO followed by treatment with a tertiary amine base. Removal of the 2xe2x80x2-acetyl group of compound 1b is readily accomplished by transesterification with methanol for 2-48 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 60xc2x0 C. to yield compound 1c. Alternative methods for deprotection of the 2xe2x80x2-acetyl group include hydrolysis in the presence of an alkali metal hydroxide or alkali metal carbonate, such as sodium hydroxide or potassium carbonate, or ammonolysis with ammonia in methanol. Compounds of formula 1d can be obtained by selective alkylation of the primary carbamate of 1c with a suitably substituted aldehyde in the presence of a reducing agent and acid. Alternatively, the corresponding acetal may be used in place of the suitably substituted aldehyde in this reaction. Preferred reagents for effecting this transformation are triethylsilane and trifluoroacetic acid in a suitable solvent, like acetonitrile, methylene chloride, or toluene at xe2x88x9220xc2x0 C. to 100xc2x0 C. Typically, the reaction is conducted for from 2-96 hours depending on the reactivity of the aldehyde or acetal. \nIt will be clear to one skilled in the art that the order of the steps in the synthetic sequence leading to compounds of the invention can be altered, provided that the functionality present in the molecule is compatible with the desired selective transformations. This is illustrated in Scheme 3. For example, compound 1a can be treated under similar conditions as described above for the reductive alkylation of compound 1c (Scheme 2) to yield compounds of the formula 1e. Removal of the 2xe2x80x2-acetyl group of compound 1e as described for the conversion of compound 1b to compound 1c (Scheme 2) provides compounds of formula 1g. Alternatively, oxidation of the 3-hydroxyl of compound 1e to the ketone of compound 1f can be conducted as described for the analogous transformation of 1a to 1b in Scheme 2. Finally, deprotection of the 2xe2x80x2-acetyl group of 1f is readily effected as described for the conversion of compound 1b to compound 1c (Scheme 2) to provide the compounds of formula 1d, wherein R is as previously defined. \nScheme 4 illustrates an alternate route for the preparation of the compounds of the invention (1d). Reaction of compound VI with a suitably substituted aldehyde and a reducing agent, such as triethylsilane, in the presence of an acid, such as trifluoroacetic acid, in a suitable solvent, such as acetonitrile, methylene chloride, or toluene, at a temperature ranging from xe2x88x9220xc2x0 C. to 100xc2x0 C. for 2-96 hours leads to the simultaneous removal of the cladinose sugar and the selective alkylation of the primary carbamate to afford compound 1e. Alternatively, the corresponding acetal of the suitably substituted aldehyde may be used to effect this transformation. Conversion of compound 1e to compound 1f and compound 1f to compound 1d can be conducted as described above. \nScheme 5, wherein R4a is halogen and R is as described above, illustrates the procedures by which compounds of formula 1b can be converted to compounds of formula 1j.\nFluorination of compound 1b can be accomplished with any one of a number of fluorinating reagents, including N-fluorobenzenesulfonimide in the presence of base, 1-(chloromethyl)-4-fluoro-1,4-diazoniabicyclo[2.2.2]octane bis[tetrafluoroborate] (SELECTFLUOR(trademark)) in the presence of base, 10% F2 in formic acid, 3,5-dichloro-1-fluoropyridinium tetrafluoroborate, 3,5-dichloro-1-fluoropyridinium triflate, (CF3SO2)2NF, N-fluoro-N-methyl-p-toluenesulfonamide in the presence of base, N-fluoropyridinium triflate, and N-fluoroperfluoropiperidine in the presence of base to give 1h wherein R4a is F. Chlorination of 1b can be effected with hexachloroethane in the presence of base, sulfuryl chloride, thionyl chloride, trifluoromethanesulfonyl chloride in the presence of base, chlorine, or sodium hypochlorite in the presence of acetic acid to give 1 h wherein R4a is Cl. Suitable brominating agents would include pyridinium hydrobromide perbromide, bromine in acetic acid, N-bromosuccinimide in the presence of base, 1,2-dibromoethane in the presence of base, or carbon tetrabromide in the presence of base to give 1h wherein R4a is Br. Suitable iodinating agents include N-iodosuccinimide in the presence of base or iodine to give 1h wherein R4a is I.\nTransformation of the halogenated derivatives 1h to the corresponding compounds of formula 1j can be accomplished through analogous synthetic routes as above. Reaction of 1h with a suitably substituted aldehyde or acetal in the presence of a reducing agent and acid yields compounds 1i. Reagent combinations for effecting this transformation include triethylsilane and trifluoroacetic acid in a suitable solvent, like acetonitrile, methylene chloride, or toluene at xe2x88x9220xc2x0 C. to 100xc2x0 C. Typically, the reaction is conducted for from 2-96 hours depending on the reactivity of the aldehyde or acetal. Compounds 1i are then converted to the corresponding compounds of formula 1j by reaction with methanol for 2-48 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 60xc2x0 C.\nIt will be understood by one skilled in the art of organic synthesis that the halogenation reaction can also be conducted at a later stage in the synthetic sequence. For example, halogenation of compound 1f (Scheme 3) affords the corresponding 2-halo derivative, which likewise can be converted to compounds of the invention by deprotection of the 2xe2x80x2-acetyl group under the previously described conditions. \nSchemes 6A and 6B illustrate the procedures by which compounds of the formula 1b can be converted to 2xcex1- and 2xcex2-fluoro derivatives of formulae 1n and 1q. Fluorination of compound 1b can be accomplished as described herein above. Reagent combinations for the conversion of compound 1b to the 2xcex1-fluoro derivative 1 k include SELECTFLUOR and sodium hexamethyldisilazide in DMF and N-fluorobenzenesulfonimide and potassium t-butoxide in THF. Typically, the reaction is conducted at xe2x88x9278xc2x0 C. to xe2x88x9260xc2x0 C. for 5 minutes to 24 hours. Reagent combinations for the conversion of compound 1b to the 2xcex2-fluoro derivative 10 include N-fluorobenzenesulfonimide and sodium hydride in DMF. Typically, this reaction is conducted at 0xc2x0 C. to 20xc2x0 C. for 1 to 24 hours.\nTransformation of the fluorinated derivatives 1k and 1o to the corresponding compounds of the invention 1n and 1q, respectively, can be accomplished through analogous synthetic routes as above. Reaction of 1k or 1o with a suitably substituted aldehyde or acetal in the presence of a reducing agent and acid yields compounds 1m and 1p, respectively. Reagent combinations for effecting this transformation include triethylsilane and trifluoroacetic acid in a suitable solvent, like acetonitrile, methylene chloride, or toluene at xe2x88x9220xc2x0 C. to 100xc2x0 C. Typically, the reaction is conducted for from 2-96 hours depending on the reactivity of the aldehyde or acetal. Compounds 1m and 1p are then converted to the corresponding compounds of the invention 1n and 1q, respectively, by reaction with methanol for 2-48 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 60xc2x0 C.\nIt will be also understood by one skilled in the art of organic synthesis that the fluorination reaction can also be conducted at a later stage in the synthetic sequence. For example, fluorination of compound 1f (Scheme 3) affords the corresponding 2-fluoro derivative, which likewise can be converted to compounds of the invention by deprotection of the 2xe2x80x2-acetyl group under the previously described conditions.\nOther compounds of the invention may also be suitable substrates for further transformation to yield other compounds of the present invention. Some of these transformations are illustrated in Schemes 7-11. \nScheme 7 illustrates the conversion of the 3-(4-nitrophenyl)-2-propenyl analog (1r) to the 3-[4-(4H-1,2,3-triazol-4-yl)phenyl]-2-propenyl analog (1t) via the intermediacy of the 3-(4-aminophenyl)-2-propenyl derivative (1s) wherein R4 is as described above. The selective reduction of the nitro group of 1r to the amine of 1s can be conducted with tin(II) chloride in ethanol at a temperature ranging from 20xc2x0 C. to 78xc2x0 C., typically for 1 to 24 hours. Alternative methods for reduction of the nitro group can also be employed, including iron/hydrochloric acid, iron/acetic acid, tin/hydrochloric acid, zinc/ammonium chloride, or sodium borohydride/nickel chloride. Conversion of the amino group of 1s to the 1,2,4-triazole of 1t is effected by condensation of the amine with N,N-dimethylformamide azine dihydrochloride in the presence of a base, such as pyridine. Reaction time is typically 2 to 72 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 115xc2x0 C. \nScheme 8 depicts the conversion of 3-(4-bromophenyl)-2-propenyl analog (1u) to the 3-(1,1xe2x80x2-biphen-4-yl)-2-propenyl analog (1v) wherein R4 is as described above. Reaction of the aryl bromide with an aryl boronic acid derivative to give the biaryl derivative is conducted under typical Suzuki coupling conditions, i.e., in the presence of a Pd0 catalyst, typically palladium tetrakistriphenylphosphine, and a base, typically sodium carbonate, potassium carbonate, potassium bicarbonate, potassium phosphate, or triethylamine in a suitable solvent, such as toluene, ethanol, methanol, DME, or THF. Reaction time is typically 2 to 48 hours at a temperature ranging from 20xc2x0 C. to 110xc2x0 C. Aryl iodides and aryl triflates are also suitable substrates for this conversion. \nScheme 9 illustrates the conversion of a substituted N-propenylcarbamate derivative (1w) to the corresponding substituted N-propylcarbamate compound (1x), wherein Rxe2x80x2 may be a member of the R group except alkenyl and alkynyl, and R and R4 is as described above. Typically, this transformation is conducted via catalytic transfer hydrogenation, in which the olefin is reacted with ammonium formate in the presence of a suitable catalyst, such as palladium on carbon, in a suitable solvent, such as methanol or ethanol, at a temperature ranging from 20xc2x0 C. to 60xc2x0 C. for 15 minutes to 24 hours. Other methods for reduction of the double bond could also be applicable, for example treatment with hydrogen in the presence of a noble metal catalyst, such as palladium or platinum. It will be obvious to one skilled in the art that the analogous N-propynylcarbamate may likewise be reduced to the corresponding N-propylcarbamate under similar conditions. \nScheme 10 illustrates a method for conversion of a secondary carbamate derivative (1j) to a tertiary carbamate derivative (1y), wherein Rxe2x80x3 is an independent member of the R group, and R and R4 are as described above, by reaction with an aldehyde and a suitable reducing agent, typically triethylsilane, in the presence of an acid, typically trifluoroacetic acid, in an appropriate solvent, such as acetonitrile, methylene chloride, or toluene. Reaction times are typically 2 to 96 hours at a temperature ranging from 20xc2x0 C. to 110xc2x0 C. Alternative methods for effecting this conversion may also be contemplated, for example reaction of a suitably protected secondary carbamate precursor with an alkyl halide in the presence of a sufficiently strong base, such as sodium hydride, potassium hexamethyidisilazide, or LDA. \nScheme 11 illustrates the conversion of compounds 1d, wherein R is as previously defined, to compounds containing a 2-hydroxy substituent (1z). This transformation may be conducted by treatment of compound 1d with charcoal in the presence of air in a suitable solvent, such as methanol or ethanol. It will be apparent to one skilled in the art of organic synthesis that other methods may be employed to effect this conversion, including for example treatment of 1d with a base, such as potassium hexamethyldisilazide or LDA, and an oxidant, such as camphorsulfonyloxaziridine or MoOPh. \nScheme 12 illustrates the conversion of 1axe2x80x2 to compounds of the invention 1bxe2x80x2, 1cxe2x80x2, 1dxe2x80x2, 1exe2x80x2, 1fxe2x80x2, and 1gxe2x80x2 with a nitrogen-containing substituent on the carbamate in the 6-position, wherein R4 and Ra are as described above. Reaction of 1axe2x80x2 with a protected aminoaldehyde derivative, such as N-(benzyloxycarbonyl)glycinal, in the presence of a suitable reducing agent, such as triethylsilane, and a suitable acid, such as trifluoroacetic acid affords the protected amine derivative (1bxe2x80x2). Typically this reaction is conducted in an appropriate solvent, such as acetonitrile, methylene chloride, or toluene, for 2 to 96 hours at a temperature ranging from 20xc2x0 C. to 110xc2x0 C. Deprotection of 1bxe2x80x2 to afford the corresponding amine derivative (1cxe2x80x2) can be readily effected by procedures known in the art, such as catalytic hydrogenation in the presence of a noble metal catalyst or catalytic transfer hydrogenation with palladium on carbon in the presence of cyclohexadiene or ammonium formate at a temperature ranging from 20xc2x0 C. to 60xc2x0 C. Typically these reactions are conducted in an inert solvent such as methanol or ethanol. Conversion of 1cxe2x80x2 to several of the compounds of the invention (Scheme 12) can be conducted with techniques known in the art, such as reductive alkylation with formaldehyde in the presence of a suitable reducing agent, such as sodium cyanoborohydride, and an acid, such as acetic acid to afford the alkylated amine derivative (1dxe2x80x2). Alternatively, 1cxe2x80x2 may be converted to the amide derivative (1exe2x80x2) by acylation with a suitably substituted acid chloride or acid anhydride in the presence of a base, such as pyridine, optionally in the presence of an acylation catalyst, such as DMAP. Amine derivative (1cxe2x80x2) may also be converted to a carbamate (1fxe2x80x2) by treatment with a suitably substituted chloroformate or pyrocarbonate derivative in the presence of pyridine. Finally, reaction of 1cxe2x80x2 with a suitably substituted isocyanate in an inert solvent, such as tetrahydrofuran or methylene chloride, provides access to the correponding urea derivatives (1gxe2x80x2). Optional deprotection of the 2xe2x80x2-acetyl group of 1bxe2x80x2, 1cxe2x80x2, 1dxe2x80x2, 1exe2x80x2, 1fxe2x80x2, and 1gxe2x80x2 is readily effected as described for the conversion of compound 1b to compound 1c (Scheme 2). \nScheme 13 illustrates the analogous procedure whereby 1axe2x80x2 is converted to compounds of the invention 1hxe2x80x2, 1ixe2x80x2, 1jxe2x80x2, and 1kxe2x80x2 with a carbonyl-containing substituent on the carbamate in the 6-position, wherein R4 and Ra are as described above. Reaction of 1axe2x80x2 with, for example, benzyl 4-oxobutanoic acid (Cannon, J. G. and Garst, J. E., J. Org. Chem. 1975, 40, 182) in the presence of a suitable reducing agent, such as triethylsilane, and a suitable acid, such as trifluoroacetic acid affords the benzyl ester derivative (1hxe2x80x2). Typically this reaction is conducted in an appropriate solvent, such as acetonitrile, methylene chloride, or toluene, for 2 to 96 hours at a temperature ranging from 20xc2x0 C. to 110xc2x0 C. Deprotection of 1hxe2x80x2 to afford the corresponding acid derivative (1ixe2x80x2) can be readily effected by procedures known in the art, such as catalytic hydrogenation in the presence palladium on carbon in an inert solvent such as methanol or ethanol. Conversion of 1ixe2x80x2 to other compounds of the invention (Scheme 13) can be executed with techniques known in the art, such as reaction with a suitably substituted alcohol or amine in the presence of a coupling agent such as dicyclohexylcarbodiimide, BOP, or PyBOP, optionally in the presence of a base, such as diisopropylethylamine, and an acylation catalyst, such as DMAP or HOAt, to afford the corresponding ester derivative (1jxe2x80x2) or amide derivative (1kxe2x80x2). Optional deprotection of the 2xe2x80x2-acetyl group of 1hxe2x80x2, 1ixe2x80x2, 1jxe2x80x2, and 1kxe2x80x2 is readily effected as described for the conversion of compound 1b to compound 1c (Scheme 2). \nA method of synthesis of 9-oxime compounds of formulae 1lxe2x80x2, 1mxe2x80x2, 1nxe2x80x2, 1oxe2x80x2, 1pxe2x80x2, 1qxe2x80x2, and 1rxe2x80x2 is illustrated in Scheme 14. Compound V is converted to compound 1lxe2x80x2 by a two-step process. Treatment of compound V with hydroxylamine in the presence of an acid such as acetic acid or formic acid in an alcoholic solvent such as methanol, ethanol, or 2-propanol at a temperature from 50xc2x0 C. to 100xc2x0 C. results in oxime formation and partial hydrolysis of the cladinose sugar. Complete hydrolysis of the cladinose to provide compound 1lxe2x80x2 is then effected by treatment with an acid, such as hydrochloric, sulfuric, chloroacetic, or trifluoroacetic, in the presence of alcohol and water. Preferred conditions for the oxime formation are the use of acetic acid in 2-propanol as solvent. Preferred conditions for cladinose removal are hydrochloric acid in aqueous ethanol. Protection of the 2xe2x80x2hydroxyl group and the oxime hydroxyl group followed by oxidation of the 3-hydroxyl provides compound 1mxe2x80x2. Protection of the 2xe2x80x2 and oxime hydroxyl groups as their acetates can be carried out by treatment of 1lxe2x80x2 with acetic anhydride in the presence of a tertiary amine base (such as triethylamine, diisopropylethylamine) or pyridine, and optionally an acylation catalyst, such as DMAP, in a suitable solvent such as methylene chloride, chloroform or THF at a temperature ranging from xe2x88x9220xc2x0 C. to 37xc2x0 C. for 2 to 48 hours. Oxidation of the 3-hydroxy group can be effected with DMSO and a carbodiimide, such as EDCI, in the presence of pyridinium trifluoroacetate in a suitable solvent, such as methylene chloride, for 1 to 24 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 37xc2x0 C. Alternative methods of oxidation include N-chlorosuccinimide and dimethylsulfide complex followed by treatment with a tertiary amine base, Dess-Martin periodinane, or oxalyl chloride/DMSO followed by treatment with a tertiary amine base. Fluorination of compound 1mxe2x80x2 with SELECTFLUOR and sodium hexamethyidisilazide in DMF provides compound 1nxe2x80x2. Compounds of formula 1oxe2x80x2 and 1pxe2x80x2 can be obtained by selective alkylation of the primary carbamate of 1mxe2x80x2 and 1nxe2x80x2, respectively, with a suitably substituted aldehyde in the presence of a reducing agent and acid. Alternatively, the corresponding acetal may be used in place of the suitably substituted aldehyde in this reaction. Preferred reagents for effecting this transformation are triethylsilane and trifluoroacetic acid in a suitable solvent, like acetonitrile, methylene chloride, or toluene at xe2x88x9220xc2x0 C. to 100xc2x0 C. Typically, the reaction is conducted for from 2-96 hours depending on the reactivity of the aldehyde or acetal. Removal of the 2xe2x80x2-acetyl and oxime acetyl groups of compounds of formulae 1oxe2x80x2 and 1pxe2x80x2 is readily accomplished by transesterification with methanol for 2-48 hours at a temperature ranging from xe2x88x9220xc2x0 C. to 60xc2x0 C. to yield compounds of formulae 1qxe2x80x2 and 1rxe2x80x2, respectively. Alternative methods for deprotection of the 2xe2x80x2-acetyl and oxime acetyl groups include hydrolysis in the presence of an alkali metal hydroxide or alkali metal carbonate, such as sodium hydroxide or potassium carbonate, or ammonolysis with ammonia in methanol. \nScheme 15 illustrates a method of synthesis of 3,9-dioxime compounds. Treatment of a 3,9-diketo compound such as a compound of formula 1d with hydroxylamine or an alkoxylamine in an alcoholic or aqueous alcoholic solvent at a temperature from 50xc2x0 C. to 120xc2x0 C. for 2 to 72 hours provides compounds of formula 1sxe2x80x2. The hydroxylamine or alkoxylamine may optionally be used in the form of an acid addition salt. The reaction may also be performed in the presence of an acid catalyst such as acetic acid, formic acid, propionic acid, or hydrochloric acid. \nScheme 16 illustrates an additional method for the preparation of compounds of the invention of Formula 1. The method is illustrated for the preparation of compounds of formula 1 uxe2x80x2 from the compound of formula 1c via the compounds of formula 1txe2x80x2 as intermediates. Compounds of formula 1txe2x80x2 (where Rd and Re are independently selected from the group consisting of hydrogen, CN, nitro, xe2x80x94C(O)Rh, xe2x80x94C(O)ORh, xe2x80x94C(O)NRhRi, xe2x80x94SO2Rh, optionally substituted C1-C8-alkyl, optionally substituted aryl, and optionally substituted heteroaryl, where Rh and Ri are independently selected from the group consisting of hydrogen, alkyl, aryl, and heteroaryl) can be obtained by reaction of 1c with a suitably substituted 1,4-dialdehyde or its equivalent in the presence of an acid. Equivalents of 1,4-dialdehydes include 2,5-dialkoxytetrahydrofurans, 1,4-dialdehyde monoacetals, and 1,4-dialdehyde diacetals. A preferred acid for effecting this transformation is trifluoroacetic acid in a suitable solvent, like acetonitrile, methylene chloride, or toluene at xe2x88x9220xc2x0 C. to 100xc2x0 C. Typically, the reaction is conducted for from 2-96 hours.\nPreferred 1,4-dialdehydes or their equivalents include 2-formyl-4,4-dimethoxybutanenitrile, tetrahydro-2,5-dimethoxy-3-furancarboxaldehyde, tetrahydro-2,5-dimethoxy-3-furancarboxylic acid methyl ester, and tetrahydro-2,5-dimethoxy-3-furancarboxylic acid ethyl ester. Compounds of formula 1txe2x80x2 can be converted to compounds of formula 1uxe2x80x2 by displacement of the pyrrole with primary amines or cyclic secondary amines in a suitable solvent, such as acetonitrile, dimethylformamide, dimethyl sulfoxide, or tetrahydrofuran, at a temperature ranging from xe2x88x9220xc2x0 C. to 120xc2x0 C. for 0.5 to 72 hours. Preferred substrates for this conversion are those in which the pyrrole is substituted with electron-withdrawing groups including, but not limited to, cyano, formyl, and alkoxycarbonyl. A particularly preferred substrate is compound 1txe2x80x2, where Rd=CN and Re=H. \nIt will be clear to one skilled in the art that the order of the steps in the synthetic sequence leading to compounds of the invention can be altered, provided that the functionality present in the molecule is compatible with the desired selective transformations. This is illustrated in Scheme 17. For example, compound 1a can be converted to compounds of formula 1vxe2x80x2 under similar conditions as described above for the conversion of compound 1c to compound 1txe2x80x2 (Scheme 16). Removal of the 2xe2x80x2-acetyl group of compounds of formula 1vxe2x80x2 as described for the conversion of compound 1b to compound 1c (Scheme 2) provides compounds of formula 1wxe2x80x2. Compounds of formula 1wxe2x80x2 may then be converted to compounds of formula 1xxe2x80x2 by treatment with primary amines or cyclic secondary amines as described for the conversion of compounds of formula 1txe2x80x2 to compounds of formula 1uxe2x80x2 (Scheme 16). Alternatively, oxidation of the 3-hydroxyl of compounds of formula 1vxe2x80x2 to the ketone of compounds of formula 1yxe2x80x2 can be conducted as described for the analogous transformation of 1a to 1b in Scheme 2. Deprotection of the 2xe2x80x2-acetyl group of compounds of formula 1yxe2x80x2 is readily effected as described for the conversion of compound 1b to compound 1c (Scheme 2) to provide the compounds of formula 1txe2x80x2. Compounds of formula 1txe2x80x2 may then be converted to compounds of formula 1uxe2x80x2 as described above in Scheme 16. \nScheme 18 illustrates a method for the synthesis of compounds of formula 1zxe2x80x2 in which Rj and Rk taken together with the atoms to which they are attached form a nitrogen-containing heterocycle. Treatment of 1c with a dialdehyde or dialdehyde equivalent in the presence of a reducing agent and acid provides compounds of formula 1zxe2x80x2. Dialdehyde equivalents include dialdehyde monoacetals, dialdehyde diacetals, 2,5-dialkoxytetrahydrofurans, 2,6-dialkoxytetrahydropyrans, and 2-alkoxy-3,4-dihydro-2H-pyrans. Preferred reagents for effecting this transformation are triethylsilane and trifluoroacetic acid in a suitable solvent, like acetonitrile, methylene chloride, or toluene at xe2x88x9220xc2x0 C. to 100xc2x0 C. Typically, the reaction is conducted for from 2-96 hours depending on the reactivity of the dialdehyde or dialdehyde equivalent. It will be clear to one skilled in the art that the reaction illustrated in Scheme 18 can also be conducted at different stages in the synthetic sequence to access compounds analogous to compounds of formula 1zxe2x80x2.\nWhen the aldehydes or acetals used in the preparation of compounds 1d, 1e, 1i, 1m, 1p, and 1y are not commercially available, they can be obtained by conventional synthetic procedures, in accordance with literature precedent, from readily accessible starting materials using standard reagents and reaction conditions. Exemplary syntheses of several of the aldehydes used in the preparation of 1d, 1e, 1i, 1m, 1p, and 1y are presented hereinafter as reference examples.\nWhen the dialdehydes or dialdehyde equivalents used in the preparation of compounds 1txe2x80x2, 1vxe2x80x2, and 1zxe2x80x2 are not commercially available, they can be obtained by conventional synthetic procedures, in accordance with literature precedent, from readily accessible starting materials using standard reagents and reaction conditions. Exemplary syntheses of several of the dialdehydes or dialdehyde equivalents used in the preparation of compounds 1txe2x80x2, 1vxe2x80x2, and 1zxe2x80x2 are presented hereinafter as reference examples.\nWhen the amines used in the preparation of compounds 1uxe2x80x2 and 1xxe2x80x2 are not commercially available, they can be obtained by conventional synthetic procedures, in accordance with literature precedent, from readily accessible starting materials using standard reagents and reaction conditions. Exemplary syntheses of several of the amines used in the preparation of compounds 1uxe2x80x2 and 1xxe2x80x2 are presented hereinafter as reference examples.\nCompounds of the invention wherein R3 is a group other than H may be prepared by methods described in WO00/75156, which is hereby incorporated by reference.\nCompounds of the invention wherein R5 is a hydroxy protecting group other than acyl may be prepared by methods analogous to those shown in the above schemes with appropriate reagents that are either commercially available or may be made by known methods.\nCompounds of the invention wherein R6 is a group other than ethyl may be prepared beginning with modified erythromycin derivatives as starting materials as described in various publications including, but not limited to, WO99/35157, WO00/62783, WO00/63224, and WO00/63225, which are all incorporated by reference herein.\nThese compounds have antimicrobial activity against susceptible and drug resistant Gram positive and Gram negative bacteria. In particular, they are useful as broad spectrum antibacterial agents for the treatment of bacterial infections in humans and animals. These compounds are particularly activity against S. aureus, S. epidermidis, S. pneumoniae, S. pyogenes, Enterococci, Moraxelia catarrhalis and H. influenzae. These compounds are particularly useful in the treatment of community-acquired pneumonia, upper and lower respiratory tract infections, skin and soft tissue infections, meningitis, hospital-acquired lung infections, and bone and joint infections.\nMinimal inhibitory concentration (MIC) has been an indicator of in vitro antibacterial activity widely used in the art. The in vitro antimicrobial activity of the compounds was determined by the microdilution broth method following the test method from the National Committee for Clinical Laboratory Standards (NCCLS). This method is described in the NCCLS Document M7-A4, Vol. 17, No. 2, xe2x80x9cMethods for Dilution Antimicrobial Susceptibility Test for Bacteria that Grow Aerobicallyxe2x80x94Fourth Editionxe2x80x9d, which is incorporated herein by reference.\nIn this method two-fold serial dilutions of drug in cation adjusted Mueller-Hinton broth are added to wells in microdilution trays. The test organisms are prepared by adjusting the turbidity of actively growing broth cultures so that the final concentration of test organism after it is added to the wells is approximately 5xc3x97104 CFU/well.\nFollowing inoculation of the microdilution trays, the trays are incubated at 35xc2x0 C. for 16-20 hours and then read. The MIC is the lowest concentration of test compound that completely inhibits growth of the test organism. The amount of growth in the wells containing the test compound is compared with the amount of growth in the growth-control wells (no test compound) used in each tray. As set forth in Table 1, compounds of the present invention were tested against a variety of Gram positive and Gram negative pathogenic bacteria resulting in a range of activities depending on the organism tested.\nTable 1 below sets forth the biological activity (MIC, xcexcg/mL) of some compounds of the present invention.\nThis invention further provides a method of treating bacterial infections, or enhancing or potentiating the activity of other antibacterial agents, in warm-blooded animals, which comprises administering to the animals a compound of the invention alone or in admixture with another antibacterial agent in the form of a medicament according to the invention.\nWhen the compounds are employed for the above utility, they may be combined with one or more pharmaceutically acceptable carriers, e.g., solvents, diluents, and the like, and may be administered orally in such forms as tablets, capsules, dispersible powders, granules, or suspensions containing for example, from about 0.5% to 5% of suspending agent, syrups containing, for example, from about 10% to 50% of sugar, and elixirs containing, for example, from about 20% to 50% ethanol, and the like, or parenterally in the form of sterile injectable solutions or suspensions containing from about 0.5% to 5% suspending agent in an isotonic medium. These pharmaceutical preparations may contain, for example, from about 0.5% up to about 90% of the active ingredient in combination with the carrier, more usually between 5% and 60% by weight.\nCompositions for topical application may take the form of liquids, creams or gels, containing a therapeutically effective concentration of a compound of the invention admixed with a dermatologically acceptable carrier.\nIn preparing the compositions in oral dosage form, any of the usual pharmaceutical media may be employed. Solid carriers include starch, lactose, dicalcium phosphate, microcrystalline cellulose, sucrose and kaolin, while liquid carriers include sterile water, polyethylene glycols, non-ionic surfactants and edible oils such as corn, peanut and sesame oils, as are appropriate to the nature of the active ingredient and the particular form of administration desired. Adjuvants customarily employed in the preparation of pharmaceutical compositions may be advantageously included, such as flavoring agents, coloring agents, preserving agents, and antioxidants, for example, vitamin E, ascorbic acid, BHT and BHA.\nThe preferred pharmaceutical compositions from the standpoint of ease of preparation and administration are solid compositions, particularly tablets and hard-filled or liquid-filled capsules. Oral administration of the compounds is preferred. These active compounds may also be administered parenterally or intraperitoneally. Solutions or suspensions of these active compounds as a free base or pharmacological acceptable salt can be prepared in water suitably mixed with a surfactant such as hydroxypropyl-cellulose. Dispersions can also be prepared in glycerol, liquid polyethylene glycols and mixtures thereof in oils. Under ordinary conditions of storage and use, these preparations may contain a preservative to prevent the growth of microorganisms.\nThe pharmaceutical forms suitable for injectable use include sterile aqueous solutions or dispersions and sterile powders for the extemporaneous preparation of sterile injectable solutions or dispersions. In all cases, the form must be sterile and must be fluid to the extent that easy syringability exists. It must be stable under the conditions of manufacture and storage and must be preserved against the contaminating action of microorganisms such as bacteria and fungi. The carrier can be a solvent or dispersion medium containing, for example, water, ethanol, polyol (e.g., glycerol, propylene glycol and liquid polyethylene glycol), suitable mixtures thereof, and vegetable oils.\nThe effective dosage of active ingredient employed may vary depending on the particular compound employed, the mode of administration and the severity of the condition being treated. However, in general, satisfactory results are obtained when the compounds of the invention are administered at a daily dosage of from about 0.1 mg/kg to about 400 mg/kg of animal body weight, which may be given in divided doses two to four times a day, or in sustained release form. For most large mammals the total daily dosage is from about 0.07 g to 7.0 g, preferably from about 100 mg to 2000 mg. Dosage forms suitable for internal use comprise from about 100 mg to 1200 mg of the active compound in intimate admixture with a solid or liquid pharmaceutically acceptable carrier. This dosage regimen may be adjusted to provide the optimal therapeutic response. For example, several divided doses may be administered daily or the dose may be proportionally reduced as indicated by the exigencies of the therapeutic situation.\nThe production of the above-mentioned pharmaceutical compositions and medicaments is carried out by any method known in the art, for example, by mixing the active ingredients(s) with the diluent(s) to form a pharmaceutical composition (e.g. a granulate) and then forming the composition into the medicament (e.g. tablets).\nThe following examples describe in detail the chemical synthesis of representative compounds of the present invention. The procedures are illustrations, and the invention should not be construed as being limited by chemical reactions and conditions they express. No attempt has been made to optimize the yields obtained in these reactions, and it would be obvious to one skilled in the art that variations in reaction times, temperatures, solvents, and/or reagents could increase the yields."}
{"text": "1. Field of the Invention\nThe present invention relates to a xe2x80x9cchip-over-boardxe2x80x9d semiconductor assembly, a type of xe2x80x9cchip-on-boardxe2x80x9d (COB) semiconductor assembly. More specifically, the present invention relates to a method and apparatus for attaching a semiconductor device to a substrate for the subsequent wire bonding of the bond pads on the active surface of the semiconductor device to circuits of the substrate.\n2. State of the Art\nDefinitions: The following terms and acronyms will be used throughout the application and are defined as follows:\nBGAxe2x80x94Ball Grid Array: An array of minute solder balls disposed on an attachment surface of a semiconductor die wherein the solder balls are refluxed for simultaneous attachment and electrical communication of the semiconductor die to a printed circuit board. A BGA may also be disposed on the printed circuit board. Conductive polymer balls or bumps may also be employed.\nCOBxe2x80x94Chip-On-Board: The techniques used to attach semiconductor die to a printed circuit board or other suitable substrate. In this instance, COB also refers to a semiconductor die attached to a printed circuit board or substrate having an aperture therein for wire bonds to extend therethrough from the bond pads on the semiconductor die to circuits of the printed circuit board or substrate.\nFlip-Chip: A semiconductor chip or semiconductor die having a pattern or array of terminations spaced around the active surface of the chip or die for face-down mounting of the chip or die to a substrate.\nFlip-Chip Attachment: A method of attaching a semiconductor die to a substrate in which the die is inverted so that the connecting conductor pads on the face of the device are set on mirror-image pads on the substrate and bonded by solder reflux or a conductive polymer curing.\nWire Bonding: Conductive wires attached between the bond pads on the active surface of a semiconductor die and the circuits of a circuit board or lead frame to form an electrical connection therebetween.\nTABxe2x80x94Tape-Automated-Bonding. Conductive traces are formed on a dielectric film such as a polyimide (the structure also being termed a xe2x80x9cflex circuitxe2x80x9d). The film is precisely placed to electrically connect a die and a circuit board or lead frame through the traces. Multiple connections are simultaneously effected.\nGlob Top: A glob of encapsulant material (usually epoxy or silicone or a combination thereof) surrounding a semiconductor die in a COB assembly.\nPGAxe2x80x94Pin Grid Array: An array of small pins extending substantially perpendicular from the major plane of a semiconductor die, wherein the pins conform to a specific arrangement on a printed circuit board or other substrate for attachment thereto.\nSLICCxe2x80x94Slightly Larger than Integrated Circuit Carrier: An array of minute solder balls disposed on an attachment surface of a semiconductor die similar to a BGA, but having a smaller solder ball pitch and diameter than a BGA.\nState-of-the-art COB technology generally consists of three techniques for attaching semiconductor die to a substrate: flip-chip attachment, wire bonding, and TAB.\nFlip-chip attachment consists of attaching a semiconductor die, generally having a BGA, a SLICC or a PGA, to a printed circuit board. With the BGA or the SLICC, the solder ball arrangement on the semiconductor die must be aligned with the connecting bond pads on the printed circuit board such that precise connection is made. After proper alignment, the semiconductor die is bonded to the printed circuit board by reflowing the solder balls. With the PGA, the pin arrangement of the semiconductor die must be a mirror-image of the pin recesses on the printed circuit board. After insertion, the semiconductor die is generally bonded by soldering the pins into place. An underfill encapsulant is then generally disposed between the semiconductor die and the printed circuit board for environmental protection and to enhance the attachment of the die to the board.\nWire bonding, unlike flip-chip attachment, generally begins with attaching either the active surface or the back side of a semiconductor die to the surface of a printed circuit board with an appropriate adhesive, such as an epoxy. In wire bonding, a plurality of bond wires are attached, one at a time, to each bond pad on the semiconductor die and extend to a corresponding lead or trace end of a circuit on the printed circuit board. The bond wires are generally attached through one of three industry-standard wire bonding techniques: ultrasonic bondingxe2x80x94using a combination of pressure and ultrasonic vibration bursts to form a metallurgical cold weld; thermocompression bondingxe2x80x94using a combination of pressure and elevated temperature to form a weld; and thermosonic bondingxe2x80x94using a combination of pressure, elevated temperature, and ultrasonic vibration bursts. The semiconductor die may be oriented either face up or face down (with its active surface and bond pads either up or down with respect to the circuit board) for wire bonding, although face-up orientation has been more common.\nIn TAB semiconductor assemblies, ends of metal leads carried on an insulating tape, such as a polyimide, are attached to the bond pads on the semiconductor die and to corresponding lead or trace ends on the printed circuit board. An encapsulant, or plastic resin, is generally used to cover the bond wires and metal tape leads to prevent contamination, aid mechanical attachment of the assembly components, and increase long-term reliability of the electronics with reasonably low-cost materials.\nA common manner of forming the encapsulant or plastic package about a semiconductor die assembly is molding and, more specifically, transfer molding. In this process (and with specific reference to COB die assemblies), after the semiconductor die is attached to the substrate (i.e., FR-4 printed circuit board), the semiconductor die assembly is placed in a mold cavity in a transfer molding machine. The semiconductor die assembly is thereafter encapsulated in a thermosetting polymer which, when heated, reacts irreversibly to form a highly cross-linked matrix no longer capable of being re-melted. In addition, another common manner of forming encapsulants for COB assemblages is xe2x80x9cglob topxe2x80x9d polymeric encapsulation. Glob top encapsulation can be applied by dispensing suitably degassed material from a reservoir through a needle-like nozzle onto the semiconductor die assembly.\nThe thermosetting polymer of transfer molding generally is comprised of three major components: an epoxy resin, a hardener (including accelerators), and a filler material. Other additives such as flame retardants, mold release agents and colorants are also employed in relatively small amounts. Furthermore, glob top encapsulation can comprise a nonlinear thixotropic material that also includes fillers to achieve the desired degree of thixotropy.\nAs previously set forth, bonding a semiconductor die to the surface of a substrate for wire bonding chip-on-board and TAB chip-on-board semiconductor assemblies is well known in the art. However, there are problems in the bonding process thereof, and specifically, problems associated with heating the adhesive between the semiconductor die and the substrate. In particular, the substrate is heated to a temperature so that the adhesive flows and then cures to thereby bond the die to the substrate with the cured adhesive therebetween. Some of the problems associated with heating the substrate at the temperature required to flow the adhesive are generally as follows: first, thermal stress exists in the interface between the adhesive and the semiconductor die and the adhesive and the substrate; and second, undesirable outgassing of a solder mask material layering a BGA on the substrate. These problems may result in health and contamination issues, and more importantly, potential detachment of the semiconductor die.\nIn U.S. Pat. No. 5,972,739, there is illustrated a semiconductor chip attached to a printed circuit board using a sheet of B-stage thermosetting resin having a filler therein, such as particles, fibers, or nonwoven fabric, wherein the active surface, or other surfaces, of the semiconductor chip is treated with a solution of a silane-coupling agent. The bumps on the active surface of the semiconductor chip extend through holes in the B-stage thermosetting resin to contact circuits on the printed circuit board.\nIn U.S. Pat. No. 5,348,607, a mixture for affixing dice to a substrate is illustrated. The mixture includes a thermoplastic polyimide, a solvent for the polyimide, and a solvent which does not dissolve the polyimide but adds thixotropicity to the mixture. The mixture is applied to the substrate, the dice are placed thereon, and the solvents are evaporated to bond the dice to the substrate. A poor solvent for the polyimide, spread over the dice and exposed portions of die attach material, causes some polyimide to precipitate out of the solution in the exposed portions of the die attach material to form a grid that extends between the dice and prevents the dice from xe2x80x9cswimming togetherxe2x80x9d during the high temperature processing. In a solvent die attachment method, the substrate is first coated with a mixture of die attach material and then the mixture is dried. Spraying a solvent over the die attach material causes the material to soften so that the dice applied thereto may adhere. The die attach material is then dried to form the bond.\nU.S. Pat. Nos. 5,990,545 and 5,866,949 illustrates a semiconductor die in a chip-scale ball grid array having a nonpolymer layer or support structure positioned between a semiconductor die and a substrate.\nU.S. Pat. Nos. 5,977,226 and 5,914,186 illustrate a vacuum dispensable silicone composition having low outgassing properties suitable for use in chip-scale packages and a high temperature antistatic pressure-sensitive adhesive tape.\nIn U.S. Pat. Nos. 5,994,772, 5,977,632, 5,929,521, and 5,892,277, various types of semiconductor attachments for attaching and electrically connecting a semiconductor die to a substrate are illustrated.\nBased on the foregoing, it would be advantageous to develop a method for attaching a semiconductor die to a substrate at lower temperatures than that done conventionally, and thus, prevent the problems as discussed previously.\nThe present invention relates to a method and apparatus for bonding an active surface of a semiconductor die to the back surface of a substrate, the substrate having an aperture therethrough for the wire bonds between the bond pads on the active surface of the semiconductor die and the circuits on the front surface to pass through the aperture. Subsequently, the aperture is filled with an encapsulant material, as well as the portions of the wires of the wire bonds being encapsulated. The present invention is directed to a method and apparatus for bonding the active surface of the semiconductor die to the substrate using an adhesive tape therebetween. The method and apparatus includes pretreating the back surface of the substrate and/or the active surface of the semiconductor die with a wetting/adhesion promotion agent layer. The wetting/adhesion promotion agent layer is applied to at least a portion of the back surface of the substrate and/or portions of the active surface of the semiconductor die.\nIn the present invention, the wetting/adhesion promotion agent layer is directed to increase the surface tension thereof to promote adhesion of the substrate and/or semiconductor die. In addition, by increasing the surface tension, the substrate may be heated at a lower temperature in order for an adhesive on the adhesive tape to flow. In this manner, the present invention provides a decrease in thermal stress in the interface between the semiconductor die and substrate. Further, the lower temperature provides a decrease in the potential for out-gassing from the adhesive used on the tape for attaching the semiconductor die to the substrate during BGA attachment of the substrate and semiconductor die to another substrate. Accordingly, the present invention provides enhanced adhesion between a semiconductor die and a substrate than that provided conventionally.\nIn addition, the wetting/adhesion promotion agent layer in the present invention may comprise a silane-coupling agent, titanate-coupling agent, and/or a solvent. The wetting/adhesion promotion agent layer may comprise one or more layers and may be applied using a dispensing method, a brushing method, a spraying method, and/or a screen printing method.\nIn a preferred embodiment, the present invention is directed to bonding a semiconductor die to a substrate for the use of wire bonding chip-on-board semiconductor assemblies.\nAlternately, the method and apparatus of the present invention may be directed to pretreating any type of semiconductor device, bare or packaged, and any type of substrate with a wetting agent layer to enhance bonding between multiple substrates, semiconductor devices, and/or a semiconductor die to a substrate, particularly where adhesive tape is used to secure the semiconductor device to a substrate."}
{"text": "Integrated circuits are generally formed on a substrate of semiconductor material out of a variety of active and passive devices. Active devices may include such devices as complementary metal oxide semiconductor (CMOS) transistors (either NMOS transistors or PMOS transistors) and may be formed from part of the semiconductor material within the substrate along with additional structures formed in conjunction with the semiconductor material. These active devices may then be interconnected using a series of conductive and insulative layers that may be formed over the active devices.\nPassive devices may include such devices as resistors. These resistors may be formed to have a desired resistance to help control the desired path and functioning of the integrated circuit. The desired resistances for the resistors may be obtained by providing a desired amount of conductivities within the material of the resistors. For example, a precise amount of dopants may be introduced into a resistor that is made from a material such as polysilicon.\nHowever, as integrated circuits and their active and passive devices have been reduced in size over the years (and as they are continuing to be miniaturized even further in the years ahead), the precision required for the manufacturing and operation of the resistors has increased. Additionally, the materials and processes used to manufacture the resistors has not kept up with the demand for high precision resistors as desired for further miniaturization while also maintaining the ease of integration for the manufacturing process that is desired for the mass production of integrated circuits.\nCorresponding numerals and symbols in the different figures generally refer to corresponding parts unless otherwise indicated. The figures are drawn to clearly illustrate the relevant aspects of the embodiments and are not necessarily drawn to scale."}
{"text": "The conventional milling operation calls for the use of an upright milling machine and a horizontal milling machine. The upright milling machine is provided with an upright milling cutter to facilitate the longitudinal cutting of a metal workpiece. Upon completion of the process of the longitudinal cutting of the workpiece, the workpiece is transferred to the horizontal milling machine on which the horizontal cutting of the workpiece is done by a horizontal milling cutter of the horizontal milling machine. It is therefore readily apparent that the conventional milling machines are inefficient and uneconomical."}
{"text": "Currently, front-facing cameras are configured on many terminal devices such as a mobile phone, a notebook computer, and a tablet computer, and are configured to shoot a photograph or a video of a target in front of the terminal devices.\nFor example, FIG. 1 shows a camera layout of an existing common terminal device such as a mobile phone. A rear-facing camera 101 is generally placed on the top or back of a screen 102 of the terminal device, and a front-facing camera 103 is generally placed on the top of the screen 102 of the mobile phone, so that the front-facing camera 103 can shoot a photograph or a video of a target in front of the terminal device.\nHowever, in an existing camera layout, relatively large space needs to be reserved at the top of a terminal device to place a front-facing camera, which causes a limited ratio of a screen to the entire terminal device."}
{"text": "In Fiber-to-the-Premises broadband network applications optical splitters are used to split the optical signals at various points in the network. Recent network specifications call for optical splitters to be incorporated in fiber distribution hubs (FDHs) which are re-enterable outdoor enclosures. These enclosures allow easy re-entry for access to optical splitters allowing splitter ports to be utilized effectively and for additional splitter ports to be added on an incremental basis.\nIn typical applications to date, optical splitters are provided prepackaged in optical splitter module housings and provided with splitter outputs in pigtails that extend from the module. The splitter output pigtails are typically connectorized with high performance low loss simple connector (SC) and/or LC connectors. This optical splitter module, or cassette, provides protective packaging for the optical splitter components in the housing and thus provides for easy handling for otherwise fragile splitter components. This approach allows the optical splitter modules to be added incrementally to the FDH, for example, as required.\nA problem may arise due to the lack of protection and organization of the connectorized ends of the splitter output pigtails. For example, these pigtails can sometimes be left dangling in a cable trough or raceway within the enclosure. Leaving an optical component, such as a high performance connector, exposed in an open area leaves it susceptible to damage. These high performance connectors if damaged can cause delays in service connection while connectors are repaired. Leaving connectorized splitter output pigtails dangling in a cabling trough also exposes them to dirt and debris in the cabling trough. In current network deployments it may be important to maintain clean optical connectors to maximize the performance of the network.\nIn addition, fiber pigtails in the current art may be organized in a manner that is not conducive to rapid service delivery. In many cases splitters may have sixteen or thirty-two output pigtails bundled together making it difficult to find a particular pigtail. Also the bundle of loose hanging pigtails can easily become entangled and/or damaged causing further delays in service delivery. These tangles can cause congestion and, in some cases, result in bend induced loss on the pigtails, causing lower system performance.\nTo solve some of these issues a separate storage tray or enclosure has been utilized to take up slack and/or store and protect splitter output pigtail connectorized ends. However, these auxiliary devices tend to take up additional space and often obscure pigtails in an enclosure making it difficult for a linesman to locate a particular pigtail and/or connector. As a result, delays may occur in deployment of new subscriber connections depending on how much time is required to access the fiber pigtails in the enclosure. Thus, there is a need for a solution that provides convenient storage for fiber pigtails and/or connectors and does not take up additional space in the enclosure. The solution should provide direct access to and identification of fiber pigtails and/or connectors.\nFinally current methods tend to result in a disassociation of the splitter module from the splitter output pigtail end. This usually results because the pigtail, once deployed, gets lost in the midst of other pigtails in the fiber jumper trough. When subscribers are taken out of service it is desirable to disconnect the splitter output and redeploy or store it for ready redeployment. It is further desirable for administrative purposes to maintain association of splitter module to splitter output pigtails so that resources are used effectively over time.\nFDHs may benefit from devices and techniques that can be adapted to facilitate the organization of fiber pigtails and fiber pigtail terminations as well as protecting sensitive optical components when not in use. The devices and techniques should also facilitate easy access to pigtails and connectors so that subscribers may be efficiently connected and disconnected from the network. The devices and techniques should also group pigtails and connectors in a manner that allows them to be associated with a particular splitter module."}
{"text": "Web assets are requested by a variety of devices with different characteristics and with different configurations—for example, these characteristics include the screen size and/or resolution, the browser, the player software, etc. In addition, different users have different browsing priorities and preferences. A customized web asset can be delivered for each different device, configuration, and user type. However, this is difficult to maintain because of the ever evolving number of different assets that must be kept up to date."}
{"text": "As wireless telecommunication networks have evolved and expanded, there has been an ever-increasing need to manage the ways in which users are charged for their wireless device usage. With the advent of smart mobile devices that are capable of putting a much bigger burden on the telecommunications network through data access and streaming, there is a need to provide management of network usage that can control access to the network in real-time, particularly with regard to data consumption. This includes everything from always-on internet access to streaming of videos and applications like games.\nUsers may pay for their access to wireless networks in both prepaid and post-paid billing systems. The wireless network includes several logical functions to handle these systems. The vast majority of subscribers use post-paid billing which is handled by an Offline Charging System (OFCS) that transfers charging information to a Billing Domain (BD) where it is processed. Prepaid billing is handled by an Online Charging System (OCS) which may also forward charging information to the BD is a similar way.\nWith the increasing usage of network resources, service providers (SP) are actively pushing tiered data consumption limits. A tiered plan allows mobile device users to get their data at one rate up to a given quota limit, and at a different rate with a second quota limit. In some alternate cases, the SP may degrade the Quality of Service (QoS) after the first quota limit and provide unlimited data consumption with the degraded QoS till the end of the billing cycle.\nWhile the QoS gating upfront with an OCS is relatively straightforward, there are no easy alternatives for the post-paid subscribers, as OCS is not in the call-control or data delivery path for them. Additionally, even with OCS, while QoS gating can be effectively applied on individual subscribers and sessions, there is no concept of applying it on a group basis. Examples of the group-based QoS gating would include enterprise users on a business usage plan, and members of a family on a personal usage plan.\nTherefore, what is needed is a method and apparatus for QoS and gating control with an offline charging system (OFCS). There is also a need to QoS and gating control as a function of quota consumption by a group of users."}
{"text": "This invention relates generally to non-slip and non-skid surfaces and more specifically to a protective cover having a non-slip surface, and a method for making such a device, with application in many fields such as frame guards for motorcycles.\nThe invention relates to the field of non-slip surfaces. The product and process of this invention may be used in any industry in which a non-slip surface has utility. One specific industry is the motorcycle racing industry. Motorcycle drivers travel on challenging terrains which include obstacles such as sharp turns, jumps, and moguls. These conditions create varying extreme forces against the driver requiring him or her to continually maintain control of the bike and his or her body. Drivers maintain control of their cycles with their hands, legs and feet. The cyclist\"\"s hands grip the steering handles while the cyclist\"\"s feet push off against pegs, extending outward from the each side of the race bike. Cyclists maintain their stability by squeezing their legs against the sides of the cycle\"\"s frame. Motorcycles for this application are generally available on the market and manufactured by companies that include Honda, Suzuki, Yamaha and Kawasaki.\nCyclists often attempt to improve their stability by improving the grip between the cyclist\"\"s legs and the cycle\"\"s frame. Traditionally, this was accomplished by the cyclist pressing his or her legs against the motorcycle\"\"s frame. This method has several disadvantages. The cycle\"\"s frame, usually made of metal, is slippery and not easily gripped. Repeated slippage and kicking of the cyclist\"\"s legs and boots against the frame causes damage to the frame\"\"s surface in the form of deformation and removal of paint and protective coatings. Consequently, the frame\"\"s raw material is exposed to the elements resulting in a more rapid corrosion and deterioration of the motorcycle frame.\nThe prior art demonstrates that some motorcycle manufacturers and owners have attempted to solve this problem by placing a protective cover, known as a frame guard, over those areas of the frame that come into contact with the cyclist. These frame guards are ordinarily made of aluminum, carbon or plastic. While they solve the problem of protecting the frame from damage and deterioration, these frames are difficult to grip, provide little friction and hinder the ability of cyclists to maintain lower body support and control.\nIn response to this problem, many drivers alter the exterior surface of the frame guards in order to increase friction and prevent slippage. This is accomplished by applying a grip tape or bed liner coating to the exterior surface of the frame guard or directly to the motorcycle frame. Each of these methods increases friction and driver stability. However, they are very temporary solutions. Both methods quickly wear down and typically separate from the frame guard after two or less uses. Some cyclists replace the grip tape or bed liner several times during a single outing. Additionally, the appearance of partially removed grip tape and liners give the motorcycle an unkempt appearance.\nAs previously explained, the disclosed product and process provide beneficial use in any field of application where an anti-slip or anti-skid surface has application. The disclosure\"\"s focus on this invention\"\"s application to a frame guard for motorcycles is not intended as any limitation on the invention\"\"s application to these other fields. Throughout this disclosure, the invention and process may be applied to any other type of surface or object capable of bonding according to the product and process described below."}
{"text": "The present invention generally relates to fabrication methods and resulting structures for semiconductor devices. More specifically, the present invention relates to a dual channel complementary metal-oxide-semiconductor (CMOS) having common gate stacks.\nIn contemporary semiconductor device fabrication processes, a large number of semiconductor devices, such as n-type metal-oxide-semiconductors (NMOS) and p-type metal-oxide-semiconductors (PMOS), are fabricated on a single wafer. In a CMOS integrated circuit (IC), complementary and symmetrical pairs of these NMOS and CMOS transistors are used for logic functions. Complementing every NMOS with a PMOS and connecting both gates and both drains together greatly reduces power consumption and heat generation relative to other logic families. For example, in a CMOS IC a high voltage on the gates will only cause the NMOS to conduct, while a low voltage on the gates causes only the PMOS to conduct."}
{"text": "The present invention relates to energetic atomic beam generation, and particularly to electrically neutral molecular or atomic beam devices and methods for mass spectrometry and surface analysis, often called fast-atom bombardment mass spectrometry (FABMS).\nFABMS has several advantages over secondary-ion mass spectrometry (SIMS). The primary advantage is that FABMS allows the use of a liquid matrix, which simplifies sample preparation and maintains a reservoir of undamaged molecular sample species when subjected to an atomic beam under dynamic (intense particle flux) conditions. Secondly, the use of a neutral atom beam in the FABMS avoids the problem of floating the ion gun system above the accelerating voltage of the spectrometer. Finally, sample charge buildup is reduced with this neutral particle bombardment.\nMolecular SIMS is invaluable for the analysis and characterization of bulk solid surfaces, their films and molecular overlayers. The static (low particle flux) used in molecular SIMS is desirable to avoid damage to the molecular solid sample. Although conventional ion sources are capable of operation under such static conditions, attempts to charge neutralize the sample with an electron flood gun to avoid sample charge buildup are often ineffective, and they can result in bombardment-induced radiation damage and electron-stimulated desorption. The dynamic (intense particle) flux of conventional atomic beam sources quickly destroys the molecular overlayers or thin molecular film samples associated with the molecular SIMS method."}
{"text": "The vascular endothelial growth factor (VEGF) proteins and their receptors (VEGFRs) play important roles in both vasculogenesis, the development of the embryonic vasculature from early differentiating endothelial cells, angiogenesis, the process of forming new blood vessels from pre-existing ones, and lymphangiogenesis, the process of forming new lymph vessels. The platelet derived growth factor (PDGF) proteins and their receptors (PDGFRs) are involved in regulation of cell proliferation, survival and migration of several cell types.\nDysfunction of the endothelial cell regulatory system is a key feature of cancer and various diseases associated with abnormal vasculogenesis, angiogenesis and lymphangiogenesis.\nAngiogenesis occurs in embryonic development and normal tissue growth, repair, and regeneration, the female reproductive cycle, the establishment and maintenance of pregnancy, the repair of wounds and fractures. In addition to angiogenesis which takes place in the healthy individual, angiogenic events are involved in a number of pathological processes, notably tumor growth and metastasis, and other conditions in which blood vessel proliferation, especially of the microvascular system, is increased, such as diabetic retinopathy, psoriasis and arthropathies. Inhibition of angiogenesis is useful in preventing or alleviating these pathological processes or slowing progression of them.\nAlthough therapies directed to blockade of VEGF/PDGF signaling through their receptors have shown promise for inhibition of angiogenesis and tumor growth, there remains a need for new or improved compounds and therapies for the treatment of such diseases."}
{"text": "1. Technical Field\nThe present invention relates to vertical free standing adjustable open frame to removably receive and hold venetian blinds and in particularly mini-blinds in an expanded xe2x80x9cdownxe2x80x9d position secured to the frame to stabilize the slats during cleaning.\n2. Background Information\nCleaning venetian blinds is an onerous task because of the blind having a number of loosely connected slats that move easily during cleaning and due to the fact the long slender individual slats are easily bent.\nIn an attempt to ease the burden of cleaning a number of different types of blind support stands have been proposed. Most of the proposed support frames provide backing against which the slats can be pressed during wiping with a wash cloth. By way of example of known support stands for use in washing venetian blinds reference may be had to the following United States Patents: U.S. Pat. No. 3,197,797 issued Aug. 3, 1965 to U. Stanley; U.S. Pat. No. 2,996,747 issued Aug. 22, 1961 to J. Iori; U.S. Pat. No. 2,908,932 issued Oct. 20, 1959 to C. Shipp; U.S. Pat. No. 2,849,745 issued Sep. 2, 1958 to E. Madsen; U.S. Pat. No. 2,634,450 issued Apr. 14, 1953 to O Britton; U.S. Pat. No. 2,598,798 issued Jun. 3, 1952 to E. Kerr; U.S. Pat. No. 2,588,557 issued Mar. 11, 1952 to F. Morris et al.; U.S. Pat. No. 2,279,691 issued Apr. 14, 1942 to R. Long et al,; and U.S. Pat. No. 676,337 issued Jun. 11, 1901 to W. Nicholls.\nNone however, known to applicant, include a portion on the frame to counteract against forces exerted lengthwise of the slat during cleaning. Also very little, if any, attention has been given to providing a simple means of releasibly anchoring the blind to the support or to adjustably varying the frame size to suite different blind widths.\nA vertical wash stand frame removably holds a mini-blind in an upright position so that the mini-blind is in the extended xe2x80x9cdownxe2x80x9d position above the floor surface for washing. The wash stand includes an adjustment for vertical height and horizontal width for holding blinds of different dimensions. It is particularly suitable for standing in the shower or bathtub and washing mini-blinds. One preferred embodiment is collapsible and portable. The open frame for supporting a venetian blind has a pair of spaced apart elongate side members each with a face plate and an end plate and are interconnected by a cross member. The cross member adjustably positions the spacing of the side members to suit blinds of various widths. A pair of saddles are adjustably mounted on the side members for releasibly supporting the top rail of the venetian blind Pins through the frame are insertable into opposed ends of the bottom rail of the venetian blind.\nAlthough any type of standard blind having slats and a head rail could be used with the present invention, miniblinds such as described and set forth in U.S. Pat. No. 5,348,068 by Elsenheimer et al. and hereby incorporated by reference, are the focal point of the present invention. The wash stand frame includes at least two vertical longitudinal leg support members which may be extendible such as with telescoping means and at least one horizontal member connoting the two leg support members together including means for adjustable extending along the horizontal axis. Each leg is held steady by a base which may hold both legs or a pair of base members defining a pair of feet may hold each leg individually.\nA principal object of the present invention is to provide a frame for holding the slats of a venetian blind captive during cleaning of the blind and a frame which is adjustable in width.\nA further principal object of the present invention is to provide a simple anchoring mean to releasibly retain the blind to be washed on the frame.\nIt is an object of the present invention to hold a mini-blind secure and taut for cleaning.\nIt is an object to the present invention to provide a support structure which may only utilize eight pieces including two side members, a top member, a bottom member, and two brackets to hold two adjustable blind holders.\nIt is an object of the present invention to provide a top member and a bottom member which are adjustable in order to accommodate every average size blind, by moving to the desired notched opening and being locked in.\nIt is another object of the present invention to provide that the top and bottom member may be connected to the connecting side member by a hinge or means for attaching permitting folding of the members for storage.\nIn keeping with the forgoing there is provided in accordance with the present invention a variable width open frame for use in washing venetian blinds comprising a pair of elongate side members each having a face plate and an end plate engagable respectively with a face portion and an end edge portion of each slat of a venetian blind being cleaned thereon, a cross member adjustably interconnecting said side members at various selected different positions selectively to vary the width of the frame, saddle means on the frame for releasibly supporting the head rail of a venetian blind and movable pin means on said frame to releasibly anchor the bottom rail of the blind to said frame."}
{"text": "1. Field of the Invention\nThe present invention relates to a system for integration of foodstuff manipulating, filling and packing equipment operations and foodstuff information analysis and audio/video data as well as a method thereof; particularly, to a system which connects a foodstuff equipment and a remote server via network, such that the operational information on the foodstuff equipment and multimedia data on the remote server can be mutually transferred, facilitating information access during foodstuff equipment operation.\n2. Description of Related Art\nCurrently known food/beverage manipulating, filling and packing equipments mostly provide functions of controlling foodstuff quantity manipulation, container position calibration, automatic tailoring or automatic ejection after package, which have been widely applied to commercial usage because of the simplicity during operation, and in particular extensively employed in the manipulation or cover-sealing operations for foodstuff service industry.\nHowever, the presently available equipments are used merely on the applications of manipulating, filling or packing; as for the products that have been processed, the equipments have never performed any recording or statistic analysis. In case that users of the said foodstuff equipments need to do statistic analysis on the quantity or types of foodstuffs which have been manipulated, filled, packed, or else on the operation status of the foodstuff equipment, then it is inevitable to input alternatively by using other recording systems.\nFurthermore, currently used equipments less frequently provide users with equipment status messages, or information related with the foodstuffs which have been manipulated, filled or packed therein; for the products manufactured, it is even rare to provide real-time advertisement channels. The promotion approaches that the users may adopt for the manipulated, filled or packed products by the equipments are usually restricted to merely print the product-related marketing propaganda or promotion contents onto the container or the plastic membrane for sealing.\nTo expand the application range of such foodstuff equipments, it needs to contemplate a working system for combining the manipulating, filling or packing operations in the foodstuff equipment as well as the information and data integration thereof, so as to enhance the its effect in industrial application, facilitating users to acquire richer and more practical information."}
{"text": "Increasingly, information is stored and processed in large data storage systems. At a base level, these data storage systems are configured with large amounts of memory to support the processing of the large amounts of data. However, the current designs limit the amount of memory that can be used by these systems."}
{"text": "1. Technical Field\nThe present invention relates generally to computer-aided diagnosis (CADx) and, in particular, to a computer-aided diagnosis (CADx) method for aiding diagnosis of three-dimensional digital image data.\n2. Background Description\nComputer-assisted diagnosis is an important technology in many clinical applications. In current clinical practice, cancer or other diseases may be missed during a physician's un-aided examination of medical image data, in part because of the large volume of data. This is particularly a problem for screening applications, since there is generally little time to devote to the examination of each patient's data, and the entire range of the data must be examined to make sure it is free from disease. Computer analysis that is performed silently in the background can greatly aid physicians in their work.\nNew technologies that offer three-dimensional (3D) scans of the human body, such as Magnetic Resonance Imaging (MRI) and Computed Tomography (CT), offer tremendous opportunities for improved detection of disease. However, the change from two dimensions to three means that there is a much larger amount of data for the physician to examine. Thus, the assistance of computer analysis becomes even more important.\nThe computer screen can only provide a two-dimensional (2D) view of a three-dimensional data set. There are many different approaches to presenting two-dimensional representations of such data, Traditional volume rendering algorithms sometimes produce satisfactory results, but they require pre-set thresholds to determine what types of tissue will be drawn as transparent and what types will be drawn as opaque. For complex anatomical structures, the volume rendering may be too cluttered for the physician to form a mental model of the underlying three-dimensional structure of the data, and how it corresponds to normal and abnormal anatomical structure. It would be useful for a CADx system to provide a visualization that helps the physician quickly build up the model, and easily distinguish abnormal structures from normal ones.\nUnfortunately, many current CADx systems are not readily accepted by physicians, because their aid is seen as more of a distraction than a help. Many such systems present the results of the computer's diagnosis to the physician by marks, such as a red circle or arrow on the softcopy, which some physicians feel can create a bias in their interpretation of the data. Furthermore, many systems are perceived as a “black box”; that is, the physicians feel that they do not have any understanding of how the systems work and how they generate their diagnosis.\nIt would be more desirable to integrate physician's knowledge into the CADx process seamlessly, and without introducing marks onto the soft-copy that may be distracting or annoying to physicians."}
{"text": "A database system may include a relational database and multiple engines that are used to query the relational database. For example, the engines may include a structured query language (SQL) engine configured to query the relational database using standard SQL, and a calculation engine configured to query the relational database based on a calculation model. There may be benefits for using one type of engine over another type of engine. For example, the calculation model is a solution to express complex data flows and calculations within the database system. The calculation model is expressed in a format that allows non-relational operations and/or relatively more complex operations that are generally not possible with standard SQL. However, the SQL engine may be better suited for other types of operations such as join re-ordering or other relational operations. Both the calculation engine and the SQL engine are associated with their own optimizer. For example, the calculation engine's optimizer may be focused on filter push down, filter combination, attribute removal, etc., while the SQL engine's optimizer may be focused on relational optimizations.\nConventionally, the calculation models are not integrated with the execution model of the SQL optimizer. For example, if the query relates to a join involving a row store view and a calculation view, the SQL engine parses, compiles, and executes the SQL query for the row store view. After execution of the row store query, the query engine instantiates, optimizes, and then executes the calculation model for the calculation view, which is then combined with the row store view."}
{"text": "Detection system based on a substrate specific enzyme such as oxygenase and dehydrogenase is used for a measurement of a substrate such as glucose in a cell culture or a blood. In particular, the measurement system based on a dehydrogenase has highly valuable as it is hardly affected by oxygen in the atmosphere. Upon reacting the dehydrogenase with the substrate, a co-enzyme such as nicotine amide nucleotide (NAD) and nicotine amide dinucleotide phosphate (NADP) is reduced (NAD+→NADH, NADP+→NADPH). Concentration of the substrate can be determined from electric current or color intensity by transferring electron from the reduced form of the co-enzyme thus formed to an electrode or a reductive chromogenic dye via an electron mediator such as PMS (phenazine methosulfate), 1-Methoxy PMS (1-methoxy phenazine methosulfate) and PES (phenazine ethosulfate). Also, the concentration of the dehydrogenase can be determined using similar system by adding the substrate externally.\nOn the other hand, the electron mediator as mentioned above has a property to accept the electron from living cells. Various methods for measuring activities of living cells via the electron mediators that carry out electron transfer such as 1-Methoxy PMS and PES have been employed in a cell proliferation assay and cytotoxicity assay used in a drug screening in drug discovery and toxicity analysis for various substances. The method for evaluating the cytotoxicity by measuring the activity of the enzyme leaked out of the part of the dead cells into a culture medium via the electro mediators. However, the sample for the measurement has to be prepared by separating the living cells from the culture medium because the electron mediators as mentioned above also carries out the electron transfer with living cells (see F. K.-M. Chan et al., Methods Mol. Biol., 2013, 979, 65-70.). The method in which the living cells are not separated is also proposed, however, the influence of the living cells cannot be avoided, which requires to lower the reaction temperature and shorten the reaction time (see JPA 2005-523022).\nIn the cytotoxicity assay, the method in which the dehydrogenase from dead cells is assayed using diaphorase, an enzyme mediating electron transfer, instead of low molecular electron mediator to suppress the electron transfer from the living cells is proposed. However, a reagent solution is of poor solubility and not easy for handling because of the use of the enzyme."}
{"text": "Launching and recovering aircraft from ship flight decks is very challenging because of the small size of the flight deck and the movement (pitch, roll, etc.) of the ship as it travels through the water. This is especially true for smaller ships like frigates and destroyers that launch and recover rotary wing aircraft. In addition, most of these ships have superstructures, functioning has aircraft hangars, in front of the flight deck. The ship's geometry, along with the combined effect of prevailing winds and the forward motion of the ship, creates an air flow over and around the ship's superstructure resulting in an air wake behind the superstructure on the ship's flight deck. This air wake is often a turbulent and unsteady vortex on the flight deck creating difficulties with aircraft maneuverability on and around the flight deck.\nThe ship's superstructure is essentially a combination of bluff bodies. The air flow separates from the sharp edges of the hangar forming shear layers and low-speed recirculation zones leading to large spatial and temporal velocity gradients in the airflow over the flight deck. Vortical flow structures are shed from the hangar and other large-scale features on the superstructure. These vertical flow structures typically have length scales similar to a helicopter's fuselage and main rotor. Therefore, as the pilot moves the helicopter through the air wake during take-offs and landings, the highly unsteady airflow causes large fluctuations in the aerodynamic loading on the rotor and fuselage that adversely affect the helicopter's lift and thrust on take-offs and landings. The air vortex can also significantly disturb the aircraft's flight path as the aircraft responds to the influence of the ship's air wake. These factors, combined with poor visibility and water spray can significantly increase pilot workload. Consequently, the margins for pilot error can be significantly reduced, which increases the frequency and severity of accidents and crash landings.\nThe impact on flight decks and pilot workload, caused by these superstructure air wakes has not been fully appreciated in ship design. Although studies have shown benefits of placing flow control devices on ship superstructures to control the air flow around the vessel and increase aircraft maneuverability during takeoff and landing, these studies have only examined modifications applied to the rearmost edges of the ships hangar to mitigate the effect of superstructure vortices to the front of the carrier flight deck as shown in FIGS. 1 and 2. These types of flow control devices have had only a slight to moderately positive, and sometimes negative, effect on the reduction of flight deck air wake.\nFIG. 1, shows several configurations, where air flow control structures extended upward from the top, aft edge of the hangar. In these locations, the structures obstruct views of the flight deck and also interfere with radar and radio communications signals to the aircraft. In other configurations, the flow control structures extend from the ends of each side of the hanger toward the flight deck. These configurations offer little to no airflow break-up and also reduce the usable area around the flight deck. Moreover, these structures are fixed in their positions, making them permanent obstructions to the radar and radio communications signals, as well as the usable flight deck area.\nFlow control structures have also been placed along the top edge, starboard side of the superstructure, as shown in FIG. 2. These structures extend out parallel to the air flow along the side of the ship, offering little to no resistance as a streamlined airflow field passes over and around them. Therefore, air vortexes typically form on the ships flight deck with the same intensity as they would have without these flow control structures in place. Consequently, a more effective flow control device is desired."}
{"text": "Modern commercial aircraft make extensive use of computer systems to control aircraft behavior, plan and execute flights, and display information to the pilots during flight operations. FIG. 1 illustrates a flight deck 40 of a Boeing 777 aircraft having a forward instrument panel 46 and a control pedestal 45 configured in accordance with the prior art. Instruments 44 and display screens 20 are distributed over the forward instrument panel 46 and the control pedestal 45 for easy access by the pilots. The display screens 20 can include primary flight displays (PFDs) 21, an engine display 27, and three multi-function displays (MFDs) 22. The MFDs 22 can present additional aircraft flight information, including navigation displays 30, aircraft checklists, communication displays, and system status information.\nAdditional instrumentation is presented at a mode control panel (MCP) 41 positioned on a glare shield 42 of the flight deck 40, and at control and display units (CDUs) 47 positioned on the control pedestal 45. The glare shield 42 can also include a display select panel 48 having three display switches 59, one for each of the MFDs 22. A series of selector switches 60 each identify one type of available information to be displayed at the MFDs 22 (e.g., one selector switch 60 identifies the electronic checklist, another identifies the navigation display, and still another identifies the communications display).\nIn operation, the pilot first depresses one of the display switches 59 to select a particular one of the MFDs 22. Then the pilot presses one of the selector switches 60 to determine what type of information will be displayed at the-selected MFD 22. One drawback with this arrangement is that, on occasion, the pilot may press a selector switch 60 without realizing that a display switch 59 had previously been selected and that the previously selected display switch corresponds to an MFD 22 other than the one the pilot wishes to control. Accordingly, it may take additional time for the pilot to first determine that the appropriate display switch 59 must be pressed before pressing a corresponding one of the selector switches 60.\nAnother feature of the arrangement described above reference to FIG. 1 is that while the display screens 20 provide all the information the pilots require for flight operations, pilots are continually seeking additional information and additional flexibility and predictability in the way the information is presented. Accordingly, it may be desirable to provide the pilot with additional information and additional options for displaying the information."}
{"text": "This invention relates to short laser pulse generators of the type which pass the sidebands generated by a laser breakdown spark.\nIn the short history of the development of the laser art and its applications, which are only now beginning to unfold, laser-induced breakdown of gases has been an easily observed and dramatic effect. Recent experiments have shown that laser-induced breakdown of gases can be accompanied by substantial spectral broadening and self-phase modulation. Of more relevance for practical utilization is the observation that the breakdown plasma cuts off transmission of most of the incident beam in a time as short as 30 picoseconds.\nThe basic concept of short pulse generation from a laser spark is to employ a spectral filter that blocks the incident laser center wavelength but transmits a sideband produced by the sudden plasma growth. Among the types of spectral filters that have been suggested are the Michelson interferometer, the Fabry-Perot etalon and the grating monochromator. Of these only the grating monochromator can have sufficient rejection ratio to be of practical interest. Nevertheless, it is an expensive and inconvenient instrument to use and conventionally passes only one sideband yielded by the plasma.\nIt is desirable to obtain a more effective spectral filter, preferably a simpler one, for use in generation of short pulses by laser-induced breakdown in gases."}
{"text": "The surface configuration of a high-gain radio-frequency transmitting and/or receiving antenna (e.g., a large unfurlable antenna for deployment from a satellite or spacecraft in extraterrestrial space) ordinarily must conform to a precisely specified configuration, typically a parabolic configuration, in order to achieve diffraction-limited performance. However, the requirements of light-weight construction and thermal stability imposed by the constraints of typical applications in extraterrestrial space often preclude such an antenna from having the rigidity necessary to ensure that the actual surface configuration remains continuously in conformity with the specified surface configuration during an extended period of antenna operation.\nEngineers involved in antenna technology customarily refer to the surface configuration of an antenna as the antenna's \"figure\". A system for monitoring the actual surface configuration (i.e., the figure) of an antenna, and for generating correction signals to change the actual surface configuration as required to maintain conformity with a specified surface configuration, is called a figure control system.\nProposals for figure control systems for use with large space-based unfurlable antennas have been described in the following documents:\n1) C. C. Huang et al., \"Structure Alignment Sensor Feasibility Demonstration\", Lockheed Missiles & Space Company, Inc., Report No. D644951, 1978.\n2) R. S. Neiswander, \"Conceptual Design of a Surface Measurement System for Large Deployable Space Antennas\", Proceedings: Large Space Systems Technology Conference, NASA Langley Research Center, 1981.\n3) P. W. Collyer et al., \"Electro-Optical System for Remote Position Measurement in Real Time\", Proceedings: Large Space Systems Technology Conference, NASA Langley Research Center, 1981.\n4) M. Berdahl, \"Surface Measurement System Development\", Proceedings: Large Space Systems Technology Conference, NASA Langley Research Center, 1981.\n5) J. M. McLauchlan, \"Spatial High-Accuracy Position-Encoding Sensor (SHAPES) for Large Space System Control Applications\", Proceedings: Large Space Systems Technology Conference, NASA Langley Research Center, 1981.\nApplications presently contemplated for large space-based unfurlable antennas require that the figure of the antenna be accurate to within 1/50 of the wavelength of the signal being transmitted and/or received, where the signal would have a frequency as high as 100 GHz. Such applications would require a measurement accuracy of about 0.06 mm in mapping the figure of the antenna at a measurement data rate of about 50 Hz. In contrast with such requirements, the measurement accuracies achievable with the antenna figure control systems previously proposed are from about 0.5 mm to about 0.15 mm (depending upon the particular system) at a measurement data rate of only about 10 Hz."}
{"text": "Rinse aid agents are applied in automatic car washes to promote the drying of wet surfaces after a vehicle is cleaned. Such compositions typically contain a hydrophobic substance, a cationic emulsifier (e.g., a dialkyl dimethyl quaternary ammonium compound, or imidazoline quaternary surfactant), solvents (usually glycol ethers and/or alcohols) and, optionally, coemulsifiers (mostly nonionic or amphoteric surfactants such as alcohol ethoxylates, amine ethoxylates, betaines, etc.). During washing of a vehicle, the use of surface-active agents in the washing water results in the formation of a continuous, firmly adhering film of water on the vehicle surface. This film has to be removed to promote drying and to avoid the formation of spots or streaks due to the presence of salts and other impurities in the water. To achieve this, surface-active quaternary ammonium compounds and hydrophobes are added to the water in the rinsing phase. Because of the adsorption of the cationic surfactant and the hydrophobe on the paint surface, the water film is opened up, and the drops of water can then easily be removed by means of a blower.\nMany of the ingredients typically used in rinse aids have undesirable environmental and performance characteristics, such as poor biodegradability or even toxicity, unpleasant smell, flammability, and high volatile organic compound (VOC) content. In addition, these compositions are often in the form of emulsions that may go through a viscous gel phase upon dilution. This makes them difficult to apply in automated carwash operations since they may plug spray nozzles and interfere with the proper operation of metering pumps.\nUS application Ser. Nos 2005/0014672 is directed to a rinse aid additive and to rinse aid compositions containing the additive. The additive includes: an amidoamine quaternary ammonium component derived from a non-animal source; a first primary amine ethoxylate derived from a non-animal source; and a second primary amine ethoxylate derived from an animal source.\nU.S. Pat. No. 6,235,914, U.S. Pat. No. 6,376,455 and U.S. Pat. No. 6,458,343 relate to quaternary ammonium compounds and formulations thereof that are useful as cleaning compositions, antistatic compounds, paper debonders, fabric softeners, hair conditioners, skin conditioners, paper deinking and ink floatation agents, asphalt emulsion agents, corrosion inhibitor agents, ore floatation agents, pesticide emulsion agents, car drying aid sprays, drilling fluid additives, and the like.\nU.S. Pat. No. 5,703,029 is directed to a water-dilutable car drybright compositions which have a strong hydrophobicizing action and which are used in rinsing liquids at carwash installations. The disclosed compositions contain ester quats which are the reaction products of alkanolamines and fatty acids, a fatty amine coemulsifier and a glycol ether solvent.\nU.S. Pat. No. 6,255,274 relates to soil release polymers in detergents, cleaning agents and fabric softeners. The latter may contain ester quats which are the reaction products of alkanolamines and fatty acids and further quaternized with customary alkylating or hydroxyalkylating agents. Triethanolamine and methyl diethanolamine type ester quaternaries are particularly preferred.\nEP 1 323 817 B1 discloses cationic preparations for cleaning hard surfaces, that contain esterquats with acyl groups derived from unsaturated C8-C22 fatty acids with iodine values of 100 to 150 as cationic surfactants.\nEP 1 840 197 A1 relates to a composition for rinsing and drying vehicles comprising at least one addition salt of an ester amine. The comparative examples disclose compositions comprising a triethanolamine oleic acid ester quat, a hydrophobe, an ethoxylated amine coemulsifier and a butyl glycol solvent.\nU.S. Pat. No. 5,827,451 discloses isotropic oil-in water microemulsions containing an oil component of a fatty acid, fatty alcohol or ester thereof; a quaternary ammonium component; an ether component; and water. The microemulsions may be used to form hydrophobic films on hard surfaces, rendering them useful in carwashes.\nEP 0 421 146 discloses biodegradable compositions that can be used as drying agents for paint surfaces. The compositions contain a cationic surfactant, an emulsifier, a solvent and an oily component.\nU.S. Pat. No. 5,391,325 discloses compositions for rinsing and drying vehicles comprising a cationic emulsifier, a fatty acid ester of an alcohol having 1 to 5 carbon atoms and a solubilizer, which can be a glycol ether solvent or an amine oxide.\nAlthough the prior art discloses compositions that are useful as rinse aids, known compositions containing biodegradable emulsifiers usually provide emulsions of inferior stability and are less efficient as rinse aids. Therefore, there is still a need for rinse aid compositions with biodegradable components that form stable emulsions and impart hard surfaces with a durable hydrophobic coating that promotes short drying times at low use concentrations. In addition, it would be highly desirable if these compositions do not form gels during dilution."}
{"text": "1. Field of the Invention\nThis invention relates to pressure-sensitive adhesives having the capacity to absorb moisture and therefore adhere to moist body surfaces for a prolonged period of time. More particularly, this invention relates to an improvement in such adhesives which affords them a combination of high plasticity and \"wet-stick\" adhesion at acceptable levels of moisture absorption over prolonged periods of time, which renders the adhesive particularly suitable for such applications as stoma seal adhesives. This invention also relates to surgical sheet materials such as adhesive tapes and to adhesive bandages and dressings, including plasters that are particularly suitable for use in connection with ostomy appliances as well as for coverings for cuts, abrasions and the like, which comprise a flexible backing member, one of whose major surfaces has adhered thereto a coating of a pressure-sensitive adhesive of the present invention.\nVarious types of pressure-sensitive adhesives have been proposed and utilized as the adhesive component in adhesive bandages, adhesive tapes and the like. Acrylate polymers, polyolefinic polymers and compounded systems based on natural or synthetic rubber polymers have all been tried or utilized as pressure-sensitive adhesives with varying degrees of success.\nA pressure-sensitive adhesive must have certain characteristics to be useful. It must be sufficiently tacky, i.e., have sufficient \"grab\" or quick-stick,\" to adhere quickly to the surface to which it is to be adhered. It must also continue to adhere to that surface over extended periods of time. A pressure-sensitive adhesive composition should also have sufficient internal strength to prevent splitting and leaving particles of adhesive on a surface to which an article coated with the adhesive has been adhered when the article is removed. Where the pressure-sensitive adhesive is designed for application to the skin, the problems of adherence are substantially increased. Although the initial tack or stick may be good, adherence over an extended period of time for many pressure-sensitive adhesives requires relatively high shear adhesion to withstand movement of the underlying skin. Also, the adhesive must be tailored to accommodate the nature of the underlying skin surface, as where perspiration and other surface changes may occur. The problem is further complicated by the fact that any pressure-sensitive adhesive designed for application to the skin must release from the skin sufficiently readily to permit removal without skin damage. Where the adhesive is too strongly adhered to the skin and has substantial internal strength, small particles of the upper layer of skin are removed with the adhesive with resulting irritation to the skin. As a result, although many pressure-sensitive adhesives are available for various commercial uses, relatively few have been found which are suitable for articles for skin applications, particular in that many of those having desirable, high sheer adhesion, have an undesirably high resistance to removal, or peel adhesion. This problem is further complicated when the adhesive is intended in such applications as those involving ostomy appliances, where it is required to maintain adhesion for prolonged periods of time while in contact with body fluids, and at the same time to prevent leakage of such fluids therethrough.\nA colostomy is the surgical creation of a new opening for the colon on the surface of the body, while an ileostomy is the surgical creation of an opening for the ileum. In an ileostomy the entire colon, the rectum, and sometimes a small portion of the ileum, is removed. In a colostomy the rectum, and sometimes a portion of the colon, is removed.\nBoth ileostomy and colostomy operations involve the creation of an artificial opening (stoma) in the abdomen to which the distal end of the healthy intestine is attached. Generally, the stoma is placed low on the abdomen and to one side. A related enterostomy for which the stoma seal adhesive of the present invention might be useful is the urinary diversion. It involves the formation of a permanent fistula through which the ureter may discharge its contents.\nAs a result of the ostomy procedure, means must be provided for conveniently and hygenically collecting human waste material from the resulting stoma. While there are many such ostomy appliances, a common problem experiences with them, as indicated above, is the establishment and maintenance of both good skin adhesion and an adequate seal to prevent waste material coming in contact with the patient's skin. This requires, inter alia, an adhesive material that has good adhesion to both dry and wet skin, adequate moisture absorptivity and sufficient elastic modulus or plasticity to prevent either inadvertent detachment from wet skin or inordinate difficulty or discomfort to the user when the device is intentionally removed."}
{"text": "The present invention relates to an emulsion for the treatment of cellulose filaments after spinning and a method for the preparation of cellulose filaments with use of the emulsion.\nUsually, in the preparation of viscose rayon filaments by a centrifugal spinning process, the cake is treated with an emulsion, dehydrated and then dried in a tunnel drier for 70 to 100 hours. The cake should be dried at a low temperature under a high moisture to be dried uniformly throughout its inside and outside. Generally, the external part of the cake is wound densely under tension so that it is difficult to shrink during drying and tends to form internal strain when dried suddenly, while the internal part of the cake is wound under low tension so that it tends to shrink by drying and forms no internal strain. Accordingly, the quality of the filaments is largely affected by the drying temperature and the drying speed.\nAccordingly, by a rapid drying method, the surface of the cake is easily dried to form a shrinking force. However, since the surface of the cake cannot shrink freely, it is dried under tension to form internal strain. Further, the shrinkage is uneven between the inside and the outside of the cake, and resultantly a difference in fineness between them occurs. In some cases, a high frequency heating system is used to dry uniformly the cake throughout its inside and outside, but it requires disadvantageously a high manufacturing cost.\nConventionally, a nonionic emulsion or a mixture of a nonionic emulsion and an anionic emulsion has been used for the emulsion treatment of cellulose filaments after spinning, and recently with the high development of surface active agents an anionic or a nonionic emulsion or their combination has been used as the base in many cases. However, no cationic surface active agent has been used.\nAn object of the present invention is to provide an emulsion for treatment of cellulose filaments such as viscose rayon filaments and acetate filaments after spinning, which reacts easily with cellulose molecule in the amorphous area of the fiber and thus lowers its swelling degree and makes free shrinkage possible along the progress of drying.\nAnother object of the present invention is to provide a method for preparation of an excellent filament having uniform fineness, smooth surface, naturally increased twist and increased tensile strength by using such an emulsion."}
{"text": "3D APIs includes common operations to resolve a multisample anti-aliased (MSAA) texture, copy resource, initialize a texture and generate mipmaps and perform other operations. In all such operations, graphics driver logic in control of a graphics processor can program a pixel shader to sample from source texture and write to a destination as render target. For texture initialization, the source is simply written to destination. For MSAA resolve and mipmap generation the graphics processor can output a simple average of the samples color into the non-MSAA destination or the next mip level. This approach enables color compression when writing to destination render target but does not utilize any cleared status that may be set for the source texture."}
{"text": "Pentasil-zeolites are defined by their structure type and more specifically by their XRD. ZSM-5 is one commercial pentasil-zeolite product.\nAs early as 1967, Argauer and Landolt (U.S. Pat. No. 3,702,886) worked out parameters for the synthesis of pentasil-zeolites, particularly those relating to the following molar ratios:\nOH.sup.- /SiO.sub.2 =0.07-10 PA1 SiO.sub.2 /Al.sub.2 O.sub.3 =5-100 PA1 H.sub.2 O/SiO.sub.2 =1-240 PA1 SiO.sub.2 /Al.sub.2 O.sub.3 =20 to 60 PA1 OH.sup.- /SiO.sub.2 =0.10 to 0.20 PA1 H.sub.2 O/SiO.sub.2 =20 to 60;\nHowever, the Argauer and Landolt procedure succeeded in synthesizing a reasonably pure phase ZSM-5 zeolite only if organic amines with a structure-giving function (i.e. template function), such as tetrapropyleneammonium compounds were used.\nSubsequent to publication of the Argauer and Landolt patent, various publications have disclosed methods of conducting the synthesis of pentasil-zeolites without requiring the very expensive, toxic and easily inflammable organic amine templates. Still other subsequent publications have disclosed substitutes for these amines. In addition to their expense, toxicity and flammability, such amines are disfavored because they are subject to thermal decomposition which can destroy the zeolite structure. Further publications have disclosed modifications of the Argauer and Landolt process directed towards improving the reactivity of the SiO.sub.2 and Al.sub.2 O.sub.3 starting materials.\nFor example, German Offenlegungsschrift 34 02 842 discloses a method wherein a specially aged aluminosilicate (which is still, however, amorphous to X-rays) is used as a nucleating gel. The method avoids the use of organic amine templates.\nEP 0 111 748 discloses aluminosilicates having a zeolite structure and methods for their preparation. The zeolite synthesis is carried out without the addition of an organic compound, but in the presence of aluminum phosphate. However, this method results in zeolites which contain phosphate.\nUnder normal conditions, the synthesis of zeolites without the use of organic compounds proceeds very slowly. Accordingly, there are no large-scale methods that do not use organic templates for the preparation of high-silica aluminosilicates having a pentasil structure (see Synthesis of High-Silica Aluminosilicate Zeolites, by P. A. Jacobs and J. A. Martens, Studies in Surface Science Catalysis 33, 1987, p.143; and Zeolites as Catalysts, Sorbents and Detergent Builders, H. G. Karge and J. Weithamp, Studies in Surface Science Catalysis, 46, 1989, p.654).\nIt is known that formation of aluminosilicates, build-up of the zeolite crystal lattice from SiO.sub.4.sup.- and AlO.sub.4.sup.- tetrahedra, nucleation of zeolites, and zeolite crystal growth all take place by way of reversible reactions. These processes depend on chemical equilibria, which can shift in different directions depending on the temperature, the hydrothermal pressure relationships and the concentrations of reactants (e.g. supersaturated or unsaturated). In synthesizing crystalline aluminosilicate materials, it is desirable to achieve as complete a conversion as possible to the crystalline aluminosilicate while avoiding both secondary phases (such as cristobalite) and an amorphous phase.\nHigh temperatures (i.e., temperatures in excess of 200.degree. C.), and consequent high reaction rates, are advantageous for the formation of crystalline aluminosilicate. However, high temperatures (in excess of 200.degree. C.) increase the probability of forming secondary phases."}
{"text": "1. Technical Field\nThe present disclosure relates to diffusion devices, and more particularly, to diffusion devices for emitting more than one active material therefrom.\n2. Description of the Background\nA multitude of active material diffusion devices or diffusers exist in the marketplace. Many of such devices are passive devices that require only ambient air flow to disperse the liquid active material therein. Other devices are battery-powered or receive household power via a plug extending from the device. A cord may be coupled between the plug and the device, or the plug may be mounted directly on the device.\nVarious means for dispensing active materials from diffusion devices are also known in the art. For example, some diffusion devices include a heating element for heating an active material to promote vaporization thereof. Other diffusion devices employ a fan or blower to generate air flow to direct active material out of the diffusion device into the surrounding environment. In another type of diffusion device, active material may be emitted from the device using a bolus generator that delivers a pulse of air to eject a scent ring. Still other diffusion devices dispense active materials utilize ultrasonic means to dispense active materials therefrom.\nIn one example a diffusion device includes two heaters for dispersion of fragrances. The device includes a housing, a plug extending from the housing for insertion into an outlet, and two containers having fragrances therein and wicks extending therefrom to absorb fragrances from the containers. Each of the heaters is disposed adjacent one of the wicks to heat the respective wick to vaporize the fragrances therein. Optionally, a CPU controlled by internal software may first activate a first of the two heaters for a predetermined period of time. After the period of time expires, the CPU deactivates the first heater and thereafter activates the second heater.\nOther diffusion devices include a housing having a cavity for receiving a cartridge. The cartridge generally has a plurality of scent elements disposed on a rotatable disk. A blower is mounted in the housing to generate airflow by passing air across a scent element and out an aperture in the housing. The housing further includes rotating means that rotate the rotatable disk, thereby rotating the scent elements thereon. The device diffuses a first scent for a predetermined time period and thereafter rotates the disk to a second scent and diffuses the second scent for the predetermined time period. This process repeats itself until the last scent element is diffused for the time period and then the disk is rotated to a home position.\nPiezoelectrically actuated vibratory type liquid atomization apparatuses are described in Helf et al. U.S. Pat. No. 6,293,474, Martin et al. U.S. Pat. No. 6,341,732, Tomkins et al. U.S. Pat. No. 6,382,522, Martens, III et al. U.S. Pat. No. 6,450,419, Helf et al. U.S. Pat. No. 6,706,988, and Boticki et al. U.S. Pat. No. 6,843,430, all of which are assigned to the assignee of the present application and which are hereby incorporated by reference herein. These patents describe an apparatus comprising a piezoelectric actuating element coupled to a liquid atomization plate. The piezoelectric actuating element vibrates the liquid atomization plate in response to alternating electrical voltages applied to the actuating element. The vibration of the plate causes atomization of a liquid supplied to it by a liquid delivery system. An electrical circuit is provided to supply the alternating electrical voltages to conductive elements that are in electrical contact with opposite sides of the actuating element. The conductive elements may also serve to support the actuating element and the liquid atomization plate in a housing that contains the device."}
{"text": "1. Field of the Invention\nThe present invention relates to an antenna cover that protects an antenna of a plasma generating device.\n2. Description of the Related Art\nIn a plasma generating device, for example, plasma electrons are accelerated by an RF electric field generated in the vicinity of an antenna in a plasma chamber and plasma is generated. At this time, if the antenna is exposed to the plasma, the antenna is sputtered, which results in becoming the pollution cause of the plasma. For this reason, a member for protecting the antenna from the generated plasma is provided. For example, a partition plate is provided between the plasma chamber and an antenna chamber to protect the antenna.\nBecause the protection member of the antenna is exposed to the plasma, the protection member is damaged by the plasma. If the protection member is abraded under an influence of the plasma and a hole is formed in the protection member, the antenna is exposed to the plasma and the plasma is polluted. If the protection member is formed thick or the antenna or the protection member is arranged to be distant from a generation position of the plasma, abrasion of the protection member can be decreased. However, generation efficiency of the plasma is deteriorated."}
{"text": "The present invention relates to the field of electro-optical signal processors. The inherent parallelism associated with optical signal processing techniques allows such systems to process a large image just as rapidly as a small one; whereas, for digital methods, the processing time tends to increase exponentially with image size. For example, a single liquid crystal television can be driven at 60 frames per second. Thus, 60 2D Fourier transforms (Fraunhofer diffraction patterns) could be formed per second. This is much faster than current state-of-the-art digital techniques. Moreover, other spatial light modulators such as ferro-electric devices, which are more expensive, can be driven at even faster rates of 2,000 frames per second.\nCurrent optical signal processors, using spatial light modulators as programmable transparencies, process a single image at a time. Since very high resolution spatial light modulators are commercially available, often the images to be processed do not require the full TV screen. Hence the \"single image\" optical systems discussed above don't take full advantage of the inherent parallelism associated with optical processing signal processing. It is thus desirable to partition the TV screen or other SLM into separate regions so that several Fourier transforms can be optically computed in parallel to significantly increase the throughput rates for many optical signal processing applications.\nUsing lenslet arrays has been suggested for spatial multiplexing. However, using lenslet arrays suffer from serious disadvantages. Most importantly, for high resolution spatial light modulators, small diameter lenses must be used. Such small lenslets tend to be poor in quality and consequently have aberrations that significantly degrade the diffraction pattern. Moreover, even if perfect small lenslets could be manufactured, their small diameter, relative to the partitioned region to be optically processed, also causes serious degradation in the resulting diffraction patterns. In order to avoid this problem, the lenslet diameter would need to be several times larger than the partitioned region; but this severely limits the number of regions that can be partitioned and thus reduces the number of Fourier Transforms that can be computed in parallel. Secondly, each of the lenslets must be aligned and spaced very accurately to within a few microns. No such lenslet arrays exist and consequently, haven't been used to perform spatial multiplexing. Thirdly, lenslet arrays are very expensive. The present invention does not suffer from any of these drawbacks."}
{"text": "The present invention relates to an optical disk and a method of producing the same and more particularly to a stamper for molding an optical disk base highly compatible with commercially available CD (Compact Disk) players, a method of producing the stamper, a method of producing an optical disk base, a method of producing an optical disk, and an optical disk base, and an optical disk.\nIn parallel with the spread of optical disks, there is an increasing demand for the timely delivery of high quality optical disks to the market. Particularly, to enhance quantity production of optical disks, it is necessary to reduce a disk base molding cycle.\nTo produce an optical disk, a stamper formed with a transfer surface is positioned in one of a pair of mold parts forming a cavity therebetween. Molten resin is injected into the cavity and then cooled off. Subsequently, the mold parts are separated in order to remove the cooled resin. As a result, the transfer surface of the stamper is transferred to the resin, forming a recording surface.\nIt is a common practice with an optical disk to hold the mold parts at a temperature of about 200xc2x0 C. lower than the temperature of resin to be injected into the cavity. This promotes the cooling and solidification of the resin injected into the cavity. Such a mold temperature is determined by the tradeoff between transferability and an increase in the tact of a disk base molding cycle. Specifically, the mold temperature should be as low as possible for increasing the tact, but would degrade transferability if excessively low. On the other hand, a high mold temperature would enhance transferability, but would increase a period of time necessary for the resin to be cooled to a parting temperature and would thereby lower the yield of optical disks.\nJapanese Patent Laid-Open Publication Nos. 7-178774, 10-149587 and 6-259815 each propose to provide a mold or a stamper with a heat insulating ability so as to enhance both the transferability and the tact of the disk base forming cycle. Specifically, Laid-Open Publication No. 7-178774 teaches a heat insulating body removably positioned in a mold in such a manner as to face the rear of a stamper. Laid-Open Publication No. 10-149587 teaches a heat insulating ceramic layer formed on a mold in such a manner as to face the rear of a stamper. Further, Laid-Open Publication No. 6-259815 teaches a stamper whose front (transfer surface) is plated with Ni (nickel) containing 20% to 30% of polytetrafluoroethylene by electroless plating. Polytetrafluoroethylene has a grain size of 1.0 xcexcm or less. The resulting Ni film is 50 nm to 70 nm thick.\nHowever, none of the above conventional technologies can enhance both the transferability and the tact of a disk base molding cycle at a high level. Laid-Open Publication No. 6-259815 has a problem that the Ni film formed on the transfer surface of a stamper obstructs the fine patterning of the transfer surface. Laid-Open Publication No. 10-149587 has a problem that the mold itself must be redesigned or replaced, wasting existing molding equipment.\nSpin coating has customarily been used to coat a molded disk base with an organic pigment which forms a recording layer because spin coating is desirable from the easy process and low cost standpoint. While the thickness distribution of the recording layer can be controlled on the basis of coating conditions, it is difficult to control the distribution of the pigment in guide grooves. Specifically, to form the recording layer, a disk base is caused to spin such that a pigment solution sequentially spreads outward over the entire disk base due to a centrifugal force. However, the centrifugal force differs from one position to another position in the radial direction of the disk base. This, coupled with the fact that the solvent evaporates while spreading outward, causes the pigment to fill outer guide grooves more easily than inner guide grooves.\nIt follows that if the guide grooves of the disk base have a uniform configuration from the inner circumference to the outer circumference, the configuration of the guide grooves filled with the pigment differs from one position to another position in the radial direction. This scatters reflectance and tracking error and other signal characteristics and makes it difficult to produce constant quality, reliable optical disks. In addition, the resulting optical disks are not satisfactorily compatible with commercially available CD players.\nJapanese Patent Laid-Open Publication Nos. 5-198011 and 5-198012, for example, disclose implementations for correcting the above difference in configuration between the inner guide grooves and the outer guide grooves filled with the pigment. The implementations are such that the configuration (depth) of the guide grooves to be formed in a disk base or a stamper is intentionally varied beforehand. None of such implementations, however, gives consideration to the decrease in the fluidity of molten resin ascribable to temperature fall. Therefore, the implementations cannot realize desirable transferability alone when a high cycle is desired, aggravating the scattering of optical disks in signal characteristics.\nIt is therefore an object of the present invention to enhance both the transferability and the tact of a disk base molding cycle at the same time.\nIt is another object of the present invention to allow a transfer surface to be finely patterned.\nIt is yet another object of the present invention to make it needless for existing molding equipment to be redesigned or replaced.\nIt is a further object of the present invention to provide an optical disk sufficiently compatible with commercially available CD players by allowing guide grooves filled with a pigment by spin coating to have a substantially uniform configuration at any position in the radial direction.\nIn accordance with the present invention, a stamper for molding an optical disk base includes a transfer surface for molding the optical disk base, and a heat insulating material extending in parallel to, but not contacting, the transfer surface.\nAlso, in accordance with the present invention, a method of producing a stamper for molding an optical disk base includes the steps of electroforming on a photoresist master having a transfer surface pattern an Ni layer having a transfer surface to which the transfer surface pattern is transferred, forming a heat insulating layer on the Ni layer, and separating the photoresist master from the Ni layer.\nFurther, in accordance with the present invention, a method of producing an optical disk base includes the steps of injecting molten resin into a cavity formed by a pair of mold parts and accommodating a stamper having a transfer surface for molding the optical disk base and a heat insulating layer extending in parallel to, but not contacting, the transfer surface, and separating the pair of mold parts to thereby remove the resin cooled off.\nFurthermore, in accordance with the present invention, a method of producing an optical disk includes the steps of injecting molten resin into a cavity formed by a pair of mold parts and accommodating a stamper having a transfer surface for molding the optical disk base and a heat insulating layer extending in parallel to, but not contacting, the transfer surface, separating the pair of mold parts to thereby remove the resin cooled off, coating a transfer surface of the resin with a recording material to thereby form a light absorption layer, forming a reflection film on the light absorption film, and forming a protection film on the reflection film.\nMoreover, in a method of producing an optical disk base, a heat insulating material is positioned beneath a recording area formed on the surface of a stamper for molding an optical disk."}
{"text": "1. Field of the Invention\nThe present invention relates to imaging apparatuses and methods. More specifically, the invention relates to imaging apparatuses and methods for characterizing objects by emitting and receiving photons.\n2. Description of Related Art\nThree-dimensional image scanning has long been a desired tool in many technical fields. Fields such as manufacturing, military technology, and medical treatments have varying applications for this technology. Three dimensional images allow for close inspection and computer image manipulation of the imaged object. Further, objects that are scanned may be rapidly reproduced and manufactured by eliminating physical measuring and computer drafting time. Early applications of this technology were limited in effectiveness and practicality because of the extremely high number of data points that were required to be processed and the speed of computers to process the data.\nThese previous applications of three-dimensional imaging apparatuses and methods also employed various inefficient methods of collecting the data points of the object. Some of these methods included physically measuring a large number of points on the object with a stylus or other device to obtain the three axes coordinates for points along an image model. Other imaging devices bombarded the object with a large number of waves or particles.\nThe waves or particles were emitted at the target and collected as they return to the source. While these procedures allowed for a large number of data points to be taken along the object, the shear volume of data points collected in the imaging processes overwhelmed the accompanying computer. Additionally, the large number of data points created a significant amount of noise in the image results, causing variation and inaccuracies in the image. These drawbacks add to the time required to image an object and the cost of the equipment. Resultantly, three-dimensional imaging has not been fully employed to its maximum potential in many fields because of these limitations.\nTherefore, what is needed is a three-dimensional imaging device and method that is capable of minimizing the number of data points required to image a three-dimensional object and that is capable of minimizing the analyzing time of the collected data.\nThe present invention has been developed in response to the current state of the art, and in particular, in response to these and other problems and needs that have not been fully or completely solved by currently available three-dimensional imaging technology. The present application discloses an imaging apparatus comprising a photon source, a photon receiver, and an analyzer. The photon source is configured to emit a photon at a target. Once the photon impacts the target, the photon reflects to the photon receiver. The photon receiver is configured to identify a reflection time of the photon and an arrival position of the photon on the photon receiver. Once the reflection time and arrival position are known, the information is transferred to the analyzer. The analyzer is configured to associate the arrival position with a position on an X-Y plane and to associate the reflection time with a Z-height of the X-Y position.\nThe imaging apparatus is capable of imaging the characteristics of two-dimensional objects, three-dimensional objects, and geographical features among other objects. The reflection time of the photon is determined from a voltage analog established by charging a voltage storage device during the interval between the photon emitting from the photon source and the photon being received by the photon receiver. In one embodiment the charging voltage is linear; in another embodiment the voltage charge is exponential.\nAn alternative embodiment of the imaging apparatus may comprise a photon source that emits a photon pulse at a target. The photons reflecting from the target impact on a photon receiver having an array of geometrically arranged pixels. The geometrically arranged array of pixels is configured to identify a plurality of photon arrival positions on the array. The imaging apparatus also has a plurality of timing circuits coupled to the geometrically arranged pixels. The timing circuits are configured to determine the reflection time of individual photons received by the photon receiver. The arrival position and reflection time are then transferred to an analyzer where the reflection times and the photon arrival position are analyzed to determine the geometric characteristics of the target.\nIn operational terms, the apparatus emits a photon pulse at a target and a plurality of timers are started in response to the emitted photon pulse. Next the photon pulse impacts the target and a number of photons reflect off of the target. The reflected photons are received on a geometrically arranged array. The locations on the array where the photons are received, establishes a plurality of arrival positions. A number of timers are also stopped upon receipt of an emitted photon on the geometrically arranged array. The time between the timer starting and stopping establishes a reflection time for individual photons. Finally, the arrival positions of the photons are compiled to generate an outline of the target in the X-Y plane and the reflection times of the individual photons are associated with positions on the outline to define the Z-components of the target. Thus, the arrival positions define the X and Y positions of the three-dimensional image and the reflection times define the Z positions of the three-dimensional image."}
{"text": "A semiconductor laser element is specified in Japanese Patent Publication No. 2002-299763 A.\nOn account of their compactness and cost-effective production, semiconductor laser diodes are used in numerous areas of application, such as, for example, data transmission, data storage, projection, material processing, optical pumping, biosensor technology and the like. In particular semiconductor laser diodes having a ridge waveguide, so-called ridge laser diodes, are used in this case. Conventional ridge technology is encountering its limits here in the face of constantly increasing requirements regarding higher powers, improved mode behavior and the like, with regard to robustness, reliability and high radiation efficiency. Central problems here are the inadequate mechanical stability, the deficient electrical loading capacity and the unsatisfactory leakage current and aging behavior of the conventional ridge laser diodes. In particular, the manufacturing yield for mass market applications and the production costs associated therewith are not optimal.\nFIGS. 6A and 6B illustrate two exemplary embodiments of a conventional ridge waveguide semiconductor laser diode in cross section. On a GaN substrate 100 there are arranged an n-conducting cladding layer 101, an n-conducting waveguide layer 102, an active semiconductor layer 2c, a p-conducting waveguide layer 103, a p-conducting cladding layer 104 and an ohmic contact layer 105. The p-conducting cladding layer 104 and the p-conducting waveguide layer 103 are etched in such a way that a ridge waveguide 200 is formed. The semiconductor laser diode thus has a main body 201 and a ridge waveguide 200.\nSide faces of the ridge waveguide and of the main body are provided with a passivation layer 107. In this case, the passivation layer 107 is arranged uniformly on the side faces, such that this is embodied in a stepped fashion. In addition, the passivation layer has an opening on the surface of the ridge waveguide, such that an electrical contact connection is made possible at this surface of the ridge waveguide.\nFor making electrical contact with the semiconductor laser diode, an n-conducting connection layer 106 is arranged on that side of the substrate which faces away from the semiconductor layers and a p-conducting connection layer 5 is arranged on the passivation layer and the ridge waveguide.\nHowever, laser diodes of this type have a deficient robustness since the ridge waveguide of the laser diode is susceptible to mechanical damage such as, for example, scratching, bond damage or external mechanical force effects. In addition, on account of the vertically formed side face of the ridge waveguide, uniform application of the p-side contact is difficult to carry out since, on account of shading effects, the contact feed at the base of the ridge waveguide is formed in a very thin fashion. However, these thinned contact regions constitute weak points in the current-carrying capacity and can lead to an electrical failure of the semiconductor laser diode. A further problem of the conventional laser diodes is a short-circuit risk. In particular on account of the passivation layer 107 having a thickness of only a few 100 μm, there is the risk of a short circuit in the region of the active layer 2c in the case where the etching process carried out for producing the ridge waveguide is implemented partially deeper, such that etching is effected near or even through the active layer 2c. "}
{"text": "Airbags for motor vehicles are known and have been used for a substantial period of time. These devices are installed on the driver and passenger side of automobiles and, in the event of a collision, are rapidly inflated with gas, to act as an energy absorbing barrier between the driver or passenger and the steering wheel or dashboard of the automobile.\nMore recently, airbag safety restraints in the form of an inflatable restraining curtain disposed along the side of a vehicle between an occupant and window or door openings have played a well recognized role in preventing injury to the occupant during a collision event. Typically, such curtains are inflated rapidly by the pressure of a reaction gas released from an inflator at the outset of generating agent in the inflator induces a chemical reaction activated by a collision signal from a collision detecting sensor when deceleration of the vehicle exceeds a certain level. The gas which is generated by the inflator is then conveyed to the airbag curtain. Inflatable curtains are typically deployed downwardly from a storage position along the roof rail so as to at least partially cover window and/or door openings across the side of the vehicle. The deployed curtain thus provides both a degree of cushioning restraint as well as a barrier preventing the occupant from being ejected from the vehicle. Due to the extended duration of a roll-over collision event where the vehicle may turn over several times it is desirable for the curtain-type airbags to remain inflated for an extended period of time so as to maintain a degree of head protection and barrier restraint until the entire event is concluded.\nCurtain-type airbags are typically formed by sewing together panels of fabric to form a plurality of pillows when inflated or by weaving two layers of fabric that are interconnected at certain areas in a Jacquard loom. Such pillowed fabrics typically utilize seams which control the shape and size of the inflated cushion. Upon inflation of these airbag cushions in response to a collision event, pressure is applied in great force, particularly on the seams and can result in seam combing. In this regard it is to be understood that the term “seam combing” refers to the phenomenon wherein applied pressure causes the seaming threads to spread apart and thereby release additional gas pressure.\nPerformance can be improved by applying substantial quantities of permeability blocking coating materials. In the past, coatings have been applied to fabrics intended for use in automotive airbags in order to resist the unwanted permeation of gas through the fabric and, to a lesser extent, to protect the fabric from detriment by the hot gases used to inflate the bags. Polychloroprene was utilized in the early development of airbags. More recently, silicone (polydimethylsiloxane or similar materials) has gained increasing acceptance.\nSince the airbag must retain its integrity during a collision event, in order to sufficiently protect the driver or passenger, there is a great need to provide coatings which provide both effective permeability characteristics and sufficient restriction of seam combing for the airbag to function optimally, if and when necessary. Thus, there exists a need for providing good adhesion and a strong bond between the individual yarns in order to effectuate long-term rigidity of the fibers to prevent unraveling at cut edges or at seams while simultaneously providing aging stability and excellent low air permeability characteristics.\nOne approach to decreasing coating weights while maintaining low permeability performance of the airbag fabrics has been to use a two layered coating system, as disclosed for example in commonly assigned U.S. Pat. Nos. 6,177,365 and 6,177,366 to Li (all hereby incorporated by reference in their entirety). Alternative coating compositions have been disclosed based on polyurethanes, such as in U.S. Pat. No. 5,110,666 or on polyurethane/polyacrylate dispersions as found in commonly assigned U.S. Pat. No. 6,169,043 (all hereby incorporated by reference in their entirety).\nCurrent industry practice typically utilizes conventional coating methods employing solvent or water based coatings for airbag fabrics. Consequently, the solvent or the water needs to be removed by drying and the organic polymer needs to be cured by subsequent exposure to heat. The coating method is predominantly some form of knife over air/gap or transfer coating.\nAn alternative approach is to laminate polymeric films to the airbag fabric. Traditional lamination methodology employs a tie coat of an adhesive layer on the fabric before film lamination. This type of approach is taught in U.S. Pat. Nos. 6,770,578, 6,908,528, and 6,753,275 (all incorporated by reference in their entirety), all of which teach application of an adhesive layer as an aqueous or solvent containing dispersion which necessitates the removal of water or solvent at a subsequent point in the process."}
{"text": "During the course of treating a fluid, for instance to control the atmospheric emission of polluting contaminants, it is common to disperse a treating agent into the fluid in order to combine with a catalyst to treat the undesired component of the fluid. In general, a mixing apparatus is used to disperse the treating agent into the fluid. Holes, which include hardened holes, holes, nozzles, injection nozzles, and the like, can be used for this purpose. Since simple ejection of the agent from such holes is not very effective in thoroughly mixing the agent with the fluid, it is also known to use baffles (sometimes also referred to as “deflectors” and “flow mixers”), usually positioned directly downstream of the injection point to encourage turbulence, thereby enhancing the mixing of the agent with the fluid.\nOne of the applications of the mixing apparatus is the treatment of flue gas containing nitric oxide and nitric dioxide (collectively, “NOx”) and other hazardous chemicals produced during the combustion of fossil fuels. NOx emissions may be minimized by catalytically reducing flue gas NOx to nitrogen and water using ammonia (NH3) in a chemical reaction before releasing the gas into the atmosphere.\nFlue gas may be treated by passing the gas through treatment systems, which attempt to provide maximum exposure of the catalyst to the flue gas in order to ensure that all the flue gas comes sufficiently into contact with the catalyst for treatment.\nIn the treatment systems, a treating agent, such as ammonia, is typically introduced into the flue gas stream by the mixing apparatus comprised of pipes having a plurality of holes. The mixing apparatus is designed to provide an even distribution of ammonia throughout the flue gas. The holes on the mixing apparatus are usually arranged such that the ammonia is injected into the flue gas towards a catalyst bed located downstream in the treatment system.\nAmmonia is commonly injected through the mixing apparatus into the flow of flue gas by utilizing an external ammonia vaporization system wherein liquid ammonia, either in an anhydrous or aqueous state, is vaporized in a heater or vaporizer, mixed with dilution air, and then routed to the mixing apparatus for injection into the flow of flue gas at a location upstream of the mixing apparatus. Typically, the ammonia is diluted with water prior to being injected through the mixing apparatus into the flow of flue gas.\nTo increase the mixing efficiency, some mixing apparatus include baffles disposed adjacent to the holes of the pipes or baffles installed between the pipes. Known mixing apparatus commonly utilize horizontal baffles for creating a turbulent mixing effect. Other mixing apparatus utilize square pipes rather than conventional circular pipes. Other known apparatus utilize baffles located in an upstream position from the pipes to create a wake downstream to increase the ammonia flue gas mixing efficiency.\nIt is therefore desirable to provide a mixing apparatus that evenly distributes the ammonia into the flue gas before reaching downstream catalyst for heat recovery, steam generation systems, packaged boilers, simple or combined cycle catalyst systems, and fired heaters."}
{"text": "Programmable logic devices (PLDs) are a well-known type of integrated circuit that can be programmed to perform specified logic functions. One type of PLD, the field programmable gate array (FPGA), typically includes an array of programmable tiles. These programmable tiles can include, for example, input/output blocks (IOBs), configurable logic blocks (CLBs), dedicated random access memory blocks (BRAM), multipliers, digital signal processing blocks (DSPs), processors, clock managers, delay lock loops (DLLs), and so forth.\nEach programmable tile typically includes both programmable interconnect and programmable logic. The programmable interconnect typically includes a large number of interconnect lines of varying lengths interconnected by programmable interconnect points (PIPs). The programmable logic implements the logic of a user design using programmable elements that can include, for example, function generators, registers, arithmetic logic, and so forth.\nThe programmable interconnect and programmable logic are typically programmed by loading a stream of configuration data into internal configuration memory cells that define how the programmable elements are configured. The configuration data can be read from memory (e.g., from an external PROM) or written into the FPGA by an external device. The collective states of the individual memory cells then determine the function of the FPGA.\nAnother type of PLD is the Complex Programmable Logic Device, or CPLD. A CPLD includes two or more “function blocks” connected together and to input/output (I/O) resources by an interconnect switch matrix. Each function block of the CPLD includes a two-level AND/OR structure similar to those used in Programmable Logic Arrays (PLAs) and Programmable Array Logic (PAL) devices. In CPLDs, configuration data is typically stored on-chip in non-volatile memory. In some CPLDs, configuration data is stored on-chip in non-volatile memory, then downloaded to volatile memory as part of an initial configuration sequence.\nFor all of these programmable logic devices (PLDs), the functionality of the device is controlled by data bits provided to the device for that purpose. The data bits can be stored in volatile memory (e.g., static memory cells, as in FPGAs and some CPLDs), in non-volatile memory (e.g., FLASH memory, as in some CPLDs), or in any other type of memory cell.\nOther PLDs are programmed by applying a processing layer, such as a metal layer, that programmably interconnects the various elements on the device. These PLDs are known as mask programmable devices. PLDs can also be implemented in other ways, e.g., using fuse or antifuse technology. The terms “PLD” and “programmable logic device” include but are not limited to these exemplary devices, as well as encompassing devices that are only partially programmable. For example, one type of PLD includes a combination of hard-coded transistor logic and a programmable switch fabric that programmably interconnects the hard-coded transistor logic.\nFIG. 1 is a simplified illustration of an exemplary FPGA. The FPGA of FIG. 1 includes an array of configurable logic blocks (LBs 101a-101i) and programmable input/output blocks (I/Os 102a-102d). The LBs and I/O blocks are interconnected by a programmable interconnect structure that includes a large number of interconnect lines 103 interconnected by programmable interconnect points (PIPs 104, shown as small circles in FIG. 1). PIPs are often coupled into groups (e.g., group 105) that implement multiplexer circuits selecting one of several interconnect lines to provide a signal to a destination interconnect line or logic block. As noted above, some FPGAs also include additional logic blocks with special purposes (not shown), e.g., DLLs, block RAM, and so forth.\nThe design of a PLD logic block can have a strong impact on the usefulness of the PLD as a whole. The lookup tables included in some PLD logic blocks, for example, can include features enabling a wide range of important functions, such as arithmetic functions or compare functions, or can make the implementation of these functions more efficient in area or speed. Improved interconnections within a logic block can also provide significant improvements to the functionality and/or performance of user designs implemented utilizing the logic block.\nFurther, a PLD interconnect structure can be complex and highly flexible. For example, Young et al. describe the interconnect structure of an exemplary FPGA in U.S. Pat. No. 5,914,616, issued Jun. 22, 1999 and entitled “FPGA Repeatable Interconnect Structure with Hierarchical Interconnect Lines”, which is incorporated herein by reference in its entirety. Additional flexibility is a valuable feature in a PLD interconnect structure.\nTherefore, it is desirable to provide improvements to a PLD logic block and/or interconnect structure that provide added flexibility, improved efficiency, and/or improved performance."}
{"text": "1. Field of the Invention\nAn aspect of the present invention relates to a sputtering apparatus for depositing an extremely low concentration of a metal catalyst on an amorphous silicon layer in order to crystallize the amorphous silicon, and more particularly, to a sputtering apparatus capable of minimizing non-uniformity of a metal catalyst caused by a pre-sputtering process without reducing process efficiency, and thus improving uniformity of the metal catalyst deposited on an amorphous silicon layer at an extremely low concentration.\n2. Description of the Related Art\nFlat panel display devices have replaced cathode ray tube display devices due to characteristics such as light weight, thin thickness, and so on, and typical examples thereof include liquid crystal displays (LCDs) and organic light emitting diode (OLED) display devices. In comparison with the LCDs, the OLED display devices are excellent in brightness and viewing angle characteristics, and require no backlight, so that OLED display devices can be realized as ultra thin displays.\nThe OLED display devices are classified into two types, a passive matrix type and an active matrix type, according to the driving method. The active matrix type OLED display devices include a circuit using a thin film transistor (TFT).\nThe TFT generally includes a semiconductor layer, which has a source region, a drain region and a channel region, and gate, source and drain electrodes. The semiconductor layer may be formed of polycrystalline silicon (poly-Si) or amorphous silicon (a-Si). The poly-Si has a higher electron mobility than the a-Si. Thus, the poly-Si is mainly used for the TFT of OLED display devices.\nAmong the methods of crystallizing the a-Si into the poly-Si, one uses a metal. The crystallizing method using the metal involves depositing a metal catalyst on a substrate using a process such as a sputtering process of depositing a metal layer on a substrate by applying plasma to a metal target formed of a crystallization-inducing metal such as nickel, or an atomic layer deposition (ALD) process of forming an atomic layer of a crystallization-inducing metal catalyst such as nickel on a substrate using a chemical method based on a reaction gas containing the metal catalyst, and crystallizing the a-Si using the metal catalyst as a seed, thereby providing an advantage in that the a-Si can be crystallized at a relatively low temperature in a short time.\nTypical sputtering apparatuses are designed to concentrate the plasma on the metal target and substrate using a magnetic assembly located at the rear of the metal target to uniformly deposit a thick layer in a short time. However, the crystallizing method using the metal degrades the characteristics of the TFT when the metal catalyst remains in the poly-Si formed by crystallizing the a-Si, and particularly, the TFT may be driven unstably when the metal catalyst remaining in the poly-Si is magnetized. Therefore, sputtering apparatuses based on the crystallizing method using the metal perform a deposition process while moving the metal target made of the metal catalyst without using the magnetic assembly to be able to prevent magnetization of the metal catalyst discharged from the crystallization-inducing metal such as nickel as well as to deposit the metal catalyst at an extremely low concentration.\nIn the sputtering apparatus that performs the deposition process while moving the metal target without using the magnetic assembly, when a pre-sputtering process of removing foreign materials such as NiO2 attached to a surface of the metal target prior to the deposition process is performed to be able to more uniformly and stably deposit the extremely-low-concentration metal catalyst, the metal catalyst discharged from the metal target during the pre-sputtering process is deposited on an edge of the substrate, so that the uniformity of the metal catalyst deposited on the substrate can be reduced.\nFurther, to prevent the metal catalyst discharged from the metal target during the pre-sputtering process from being deposited on the substrate, a region where the plasma is produced is isolated during the pre-sputtering process. In this case, a separate shield is required to isolate the plasma production area. Further, it takes a predetermined time to form or remove the separate shield depending on progress of the pre-sputtering process, so that the deposition process is delayed, and thus overall process efficiency is reduced."}
{"text": "This invention relates to intraocular surgical instruments and more particularly to instruments adapted to vitrectomy surgery that incorporate reciprocating cutters.\nIt is sometimes necessary to surgically remove a portion of the vitreous body of the eye in order to eliminate opaque regions of the body, blood, retinal tissue and the like. The vitreous body has a jelly-like consistency, requiring that portions to be excised be cleanly cut off from the remaining portion of the vitreous body. One method of severing vitreous portions involves the use of a surgical instrument that comprises coaxial slender tubes and which is inserted into the region of the vitreous to be removed. The outer tube has an orifice with a sharp cutting edge on the internal surface of the tube, while the inner tube has a sharp cutting edge formed on its outer periphery. The portion of the vitreous body to be excised is pulled into the orifice in the outer tube by suction applied to the inner tube, and the portion is severed by moving the cutting edge of the inner tube past the cutting edge of the outer tube in a reciprocating fashion.\nIt is important that the cutting edges of the reciprocating tubes of such instruments fit together closely in order to provide a good cutting action. If the cutting edges do not cooperate properly, there is a danger that a portion of the vitreous body might become trapped between the outer wall of the inner tube and the inner wall of the outer tube. This is a serious problem because tissue that becomes trapped is difficult to cut cleanly away from the remaining vitreous body, and further movement of the instrument may cause tension on the vitreous body that may damage the retina.\nVarious means have been provided for effecting the proper orientation of the cutting edges of the inner and outer tubes in order to clearly sever the portion of vitreous body to be removed. In most of the prior art devices, this objective is addressed by providing minimal clearance between the internal wall of the outer cutter tube and the external wall of the inner cutter tube along the entire length of their adjacency. While this arrangement helps to ensure the proper orientation of the tubes at the point where the cutting edges meet, it has caused a considerable amount of friction between the two walls.\nIn an attempt to overcome the problem of friction while retaining a proper alignment at the cutting edges, a longitudinal portion of the outer tube has been cut away between its two ends, as disclosed in U.S. Pat. No. 4,246,902, issued to Martinez. In this fashion, the friction that would have otherwise occurred between the cut-away portion and the inner tube is of course eliminated, but the amount of friction caused by movement of the remaining portions against the inner tube is maintained, and the lessened structural integrity caused by the removal may also be a source of problems. In particular, the cutting away of a large longitudinal portion of one of the cutting tubes may increase the flexibility of that tube to an extent that the precise orientation of the cutting edges is lost.\nIn addition, prior art devices of this type are generally manufactured in a way that increases the cost of the devices and causes disposability after a single use to be impractical. For example, in U.S. Pat. No. 4,246,902, the device is generally complex and requires a number of threaded parts."}
{"text": "In the prior art as presented by U.S. Pat. No. 5,546,165, by Rushing et al, there is disclosed an electrostatographic copier/printer with adjustable process control parameters wherein image contrast and color balance can be adjusted. The process control parameters considered adjustable are initial or primary charge, V.sub.o, exposure E.sub.o and developer bias V.sub.B.\nThe copier/printer described in Rushing et al is specifically directed to overcoming the problem of non-uniformities in reproductions that are associated with a direction that is cross-track to the process direction. Such non-uniformities cannot be well accommodated by making overall process changes.\nIn order to overcome this problem, the copier/printer described in Rushing et al provided for the printing of special test images. These printed test images are then delivered and fed into a scanner. The scanned data are analyzed to characterize the overall process in terms of cross-track uniformity and deviation from desired tone scale. The Rushing et al patent recognizes that cross-track nonuniformity and deviation from desired tone scale may result from misalignment of one or more process stations for example an exposure station, a charging station or a development station. Thus, z-axis skewing (tilt) of the LED printhead and/or Selfoc lens which focuses light from the LEDs (light-emitting diodes) onto the photoconductive surface can provide for a non-uniformity. Other recognized potential problems are the tilt of the primary charger or non-parallelism of the development station relative to the photoconductor. The Rushing et al patent proposes that, as a result of the scanning, adjustments can be made in LED exposure on-time or driver current to correct for same.\nA problem with the apparatus of the prior art is that a scanner is required to provide for correction for such non-uniformities. The tilt and bow non-uniformities are the most common and most objectionable.\nIn the printer art generally, it is also known to detect non-uniformities using a visual detection by an operator and correction can be provided by adjustment of signals to each of the recording elements. In this regard, reference may be had to U.S. Pat. No. 5,596,353.\nA problem with the above apparatus is that where several thousand recording elements are arranged on the printhead, it becomes relatively cumbersome to provide the manual correction suggested in this patent except where only a relatively few recording elements require correction, a situation not typical in a tilt or skew or bow situation.\nIt is an object of the invention therefore to provide for an improved control of cross-track non-uniformity in a reproduction apparatus.\nAccordingly, it is an object of the present invention to provide an electrostatographic copier and/or printer apparatus and method for adjusting for cross-track non-uniformities without the need of a scanner to scan a test print printed by the apparatus and without the need for an operator to identify specific recording elements and provide specific corrections to specific recording elements."}
{"text": "The present invention relates to catalysts for purification of gaseous condensate and oil fractions from mercaptans and can be used in oil-processing and in petrochemical industry.\nUsually oxidation of mercaptans is performed by oxygen or air at room or elevated temperature in the presence of homogenous or heterogeneous catalyst based on a transitional metal:2RSH+½02→RSSR+H20\nMercaptans are extracted from organic phase by strongly alkali aqueous solutions, in which at elevated temperature and pressure catalytical oxidation of mercaptans takes place, and as a catalyst salts from metals of alternating valency (Fe, Co, Ni, V, Mn, Cr) in the form of simple or complex compositions, more often phatocyanites (as disclosed for example in European patent no. 394571, German patent 3008284 and others). The catalysts of this type provide sufficiently complete removal of mercaptans. Their common disadvantage is not sufficient availability and high cost of phatocyanites, which leads to making the process of purification expensive. Moreover, the necessity to conduct the process in strongly-alkli medium leads to a significant corrosion of equipment and worsening of the properties of oil pipelines.\nIt is known that phtalocyanines in catalytical compositions can be replaced with other compounds of transitional metals.\nA catalyst are known based on metallo-organic derivatives of metals of VIA, VIIA, and VIII groups with general formula Mea(R)x(CO)y where R is aromatic ligand, for example benzol or its alkyl derivatives, antrazen, benzpirin, phenanthren (U.S. Pat. No. 3,053,756). Disadvantages of such catalysts include their low stability, high cost and toxicity.\nA catalyst for oxidating demercaptanization of carbohydrates is known which is a helate complex of a transitional metal with bi, tri, or tetra dentant ligand, containing at least one amide group (French patent 2,573,087). In particular, replaced 2-(alkyl(aryl, alkylary))-aminocarboxypyriditines and others even more complicated compounds are utilized. As a metal it can be Co, Fe, Cu, Ni, Mn. The disadvantage is a low stability of the catalyst and use for its manufacture of expensive and scarce components.\nAlso, catalysts of oxidating demercaptanization are known-complexes of copper with tetracyantiophenol and tetracyandentin (French patent 2591610). These catalysts provide high degree of purification, but their practical use is questionable because their very high cost and scarce availability of the components.\nA method of oxidative demercaptanization are known which is performed by oxidation of mercaptans with oxygen of air in presence of helate complexes of a transitional metal (Co, Fe, Cu, Ni, Mn) with polydentan ligand from the class of amides, in particular from aminocarboxyperidines (French patent 2573087). The main disadvantage of the method with the use of this catalyst is high cost of its components.\nA heterogeneous catalyst is described, based on complexes of copper with amino derivatives (aminoalcohols, aminoacids, amines) applied on a mineral carrier or activated coal (European patent 996500). The disadvantage of this catalyst is a low content of active phase on the surface of the carrier, which inevitably leads to a significant consumption of the heterogeneous catalyst.\nThe closest solution to the present invention is disclosed in U.S. patent no. In accordance with this patent, a catalyst of alkali-free purification of oil fractions from mercaptans is proposed, which contains a copper oxide, a complex of copper with nitrogen-containing compound of amines, aminoalcohols, aminoacids, or amids and an inert carrier. The disadvantage of this method is a complex technology for producing catalyst and a high consumption of the carrier, which makes the catalyst quite expensive."}
{"text": "This disclosure generally pertains to a music processing system that converts sound from musical instruments into an electronic data format. More specifically, this invention pertains to a system and method that converts sound generated by musical instruments to a form to be used in electronic media based on a first harmonic of an input signal. In one embodiment, the data format is the Musical Instrument Digital Interface (MIDI) format.\nFor years digital keyboard players enjoyed unparalleled flexibility and functionality in interfacing and composing with their computers, such as the ability to instantly create notation and change sounds generated by their instruments with the push of a button. The music processing system described herein offer this flexibility and functionality to guitarists as well as the ability to use a guitar with computer games. The methods and apparatus described may comprise a pick-up and converter that attaches directly to any electric, acoustic electric or acoustic guitar, thereby making a user's guitar fully plug and play compatible with Windows XP or higher as well as Mac OSX. Preferably, no driver installation is necessary.\nThe music processing system described herein may be adapted for use with Guitar Wizard, a game that allows users to jam along to popular songs while learning to play a real guitar. Guitar Wizard teaches aspiring musicians everything from single note picking to complex chords and strumming techniques. Modern Digital Audio Workstation (DAW) software, such as Sony Acid™ Music Studio and Apple GarageBand harness the power of PCs, allowing musicians to play samples and software instruments. With the music processing system described herein, guitarists can control these programs to play sampled sounds and synthesized instruments such as a keyboard or piano, a different style guitar, drums or a woodwind instrument. Using the music processing system described herein, guitarists can compose a complete masterpiece controlling and recording each instrument from trumpets to tympanis using their guitar.\nUsing the music processing system described herein, users will enjoy the ability to connect a real guitar to console systems bridging the gap between gaming and reality. For instance, using the music processing system described herein, one may be able to: use a guitar to connect with a computer, operating with for instance Windows XP and/or Mac OSX; learn to play guitar; record, compose and edit music easily; arrange with flexibility and control; and convert recorded songs into sheet music. As described below, the pick-up and control components of the music processing system mount on any guitar and preferably recognizes and transmits specific instructions for each individual note played on the guitar, thereby allowing for great flexibility in playing and recording. This is conveyed simply as a list of events which describe the specific steps that a soundcard, program or other device use to generate the specific sound. At its simplest the language would indicate for example ‘Middle C on” at a specific time along with the volume of the note—then it would indicate “Middle C off” at a later time. Any number of other commands can be added to make it as expressive as desired.\nThus, the music processing system may allow the user to make his or her guitar sound like another instrument. With the system, a guitar can sound like anything: a keyboard or piano, a completely different style guitar or a guitar with any number of different effects applied, a woodwind or brass instrument or the human voice. Each note can even be assigned to play a different recorded clip or sound effect. Different or “drop” tunings are simple because the note or tuning of the guitar need not be changed. The instructions for playing the note are simply “transposed” to the desired note in accordance with the desired tuning. The language of the music processing system is very specific as to what note is being played down to the specific fret on each string. This information can be used in conjunction with a learning program to teach guitar. Since each string is tracked individually this can be a very complex and robust application, teaching everything from single note picking to complex chords and strumming techniques.\nFurthermore, the instructions generated can optionally be recorded on a computer memory. This allows recorded instructions to be edited using computer software. A single note within a recorded song is easily adjusted because all that is changed is the instruction for that specific note. To change or delete a note or passage in a regular recording would require clipping out the undesirable portion and re-recording—not an easy task as precision is next to impossible yet required. An embodiment described herein also allows for easy tempo changes of a recorded performance. The instruction is simply adjusted to change the tempo, thereby avoiding pitch change when a recording slowed down. Editing recorded music is simple using computer software—drag and drop functionality may be provided to edit individual notes. Shorten or lengthen a note simple by clicking on it and changing its duration. Using software, the user can change the whole recording to a new key using the same principle described above regarding alternate tunings. Users can cut and paste a section for use later in the song. File size is small because the methods described herein store instructions for playing a note, not sampling and digitizing the actual note or sound wave. This saves storage space on a hard drive. For example a sampled or digitized 1 minute clip requires about 10 Megabytes of data. The same 10 seconds with the music processing system only requires 10 Kilobytes for the same 1 minute clip. Many files that are already recorded in this language have tracks that are separated from the rest of the tracks making it easy to listen to just one instrument track and study it to learn more about it or how to play it. Then this track can be muted, played over to practice playing the song or for a live performance with backing tracks. One may print out actual sheet music of what has been recorded. It is very simple to convert the recorded instructions into musical notation. One may also create a ring tone for a cell phone.\nThe music processing system can be used to trigger much more than notes. The instructions for turning a note on and off and other such commands can optionally be used to activate any action or event within a program or computer game. Many prerecorded elements such as loops or tracks can be triggered on a computer program turning them on and off as backing tracks for example. These elements can also be turned on and off building them to create a song by selecting each individual element by playing a single note. Events can optionally be triggered in games. Playing a specific note or notes can be used for a game such as Guitar Hero™ or other similar game. It can be as simple as the current offerings or as complex as a real-world guitar performance. Notes could even be used to move a character around the screen. Embodiments of the present invention can also be used in conjunction with a wide variety of musical equipment. Most electronic musical equipment supports one of the various versions of the MIDI format.\nOne using the technology disclosed herein can achieve sound electronically using any classical instruments or any sound source. One method disclosed herein includes analyzing sound from the sound source, and then generating an appropriate sound electronically based on the detection of the first harmonic of the input signal. The second method disclosed herein requires fast and precise first harmonic period determination from the signals generated by a classical instrument, and then the measured period may be transformed to digital information acceptable by electronic instruments to generate sound electronically. Algorithms for transforming the measured period into digital information are disclosed in a co-pending patent application entitled “Adaptive Triggers Method for Signal Period Measuring,” U.S. application Ser. No. 11/873,970, filed Oct. 17, 2007, the disclosure of which is incorporated by reference herein. However, other tone detection methods known in the art may also be used. Such algorithms, which for instance provide a solution for transforming guitar sounds to MIDI commands, require powerful thirty-two bit microprocessors and/or DSP processors, as will be described below."}
{"text": "This invention relates to a magneto-optical recording capable of recording medium, reproducing/erasing with a laser beam, and more particularly to the structure of a magneto-optical disk effective for the suppression of deterioration in the reproduction output due to repetitions of recording/erasing, while maintaining a high S/N and also for the control of shape of a formed recording domain.\nWith the recent progress toward a highly informational society, needs for a file memory of larger capacity and higher density are now increasing and optical disks have been regarded as a response to these needs. Particularly, magneto-optical recording is now to be used in the form of a practical, rewritable optical disk. Currently known layer structures are classified into two types, i.e. three-layer structure (underlayer/recording layer/protective layer) and four-layer structure (underlayer/recording layer/protective layer/metallic layer). Most importantly, disks of a four-layer structure utilize both the Kerr effect and the Faraday effect of magneto-optical recording layer concurrently and thus a large rotational angle can be obtained and, also the reproduction output can be increased at the driving. Therefore, disks of a four-layer structure are now extensively studied and are under development.\nOn the other hand, studies on higher density and higher speed are also underway. For example, it has been proposed to utilize a field modulation recording system, a system for using a double layer based on an exchange couple, etc. in order to attain the overwritting. In order to attain a higher density, it has been studied to form micromagnetic domains or to use pit edge recording. In any event, disks having a high S/N and a good controllability of magnetic domains are required. If magneto-optical disks can be used not only in computer file memory, but also for domestic purposes for example as rewritable compact disks, a larger requirement will be posed on the development. The first condition for this purpose is a lower cost. To satisfy the condition, it is necessary to use low cost plastics as substrate materials and to make the disk structure as simple as possible. This can be attained only as a result of of improvements in disk technology such as (1) technology of preparating substrates, (2) technology of making a magneto-optical recording layer of higher corrosion resistance and (3) optimization of disk structure, etc. In particular, the technology of making a magneto-optical recording layer of higher corrosion resistance, as mentioned in (2), is important.\nIn the above-mentioned disk technology, research and development have been thus far concentrated on an increase in the rotational angle, while no particular consideration has been paid to the decrease in disk characteristics, especially reproduction output, when recording/erasing are repeated, and therefore no higher reliability has been obtained. This decrease in disk performance results from the amorphous state of magneto-optical recording layers being relaxed by a laser beam used for the recording/erasing.\nIn the case of higher density recording by field modulation recording, recording domains themselves interfere with one another or no better shapes of recording domains are obtained, and thus no satisfactory reproduction output is obtained. This is because parts in a larger ratio of the domain wall to the domain area exist and consequently the energy of the domain wall is so high that the shapes of recording domains become unstable.\nFurthermore, in the case of domestic applications, etc. on the other hand, no satisfactory simplification of layer structure has been made yet, resulting in poor mass production and higher cost. In other words the magneto-optical disks have been found to be unsuitable for the domestic applications. Still furthermore, an underlayer made mainly of inorganic compounds must have a thickness of at least 750 .ANG., in view of the optical effect. This increases an overall stress of the magneto-optical recording medium and causes the underlayer to peel off, especially when formed on a plastic substrate. In order to prevent such a peeling, it has been attempted to reduce the stress on the underlayer of inorganic compounds or to apply some treatment to the substrate surface. The thickness of the underlayer of inorganic compound has been restricted to 700-900 .ANG. in view of the optical characteristics, which case the recording layer has been sometimes corroded via pinholes existing in the underlayer."}
{"text": "In a TDD (Time Division Duplexing) system, information is sent and received by using different time slots of a same frequency carrier. With regard to a frame structure of the TDD system, FIG. 1 may be referred to, in which each radio frame has a frame length of 10 ms and includes ten subframes each of which has a length of 1 ms, and each half radio frame has a frame length of 5 ms.\nIn the TDD system, an uplink-downlink configuration may be semi-statically adjusted according to different types of services, so as to meet a service requirement of asymmetrical uplink-downlink. In a 3GPP (3rd Generation Partnership Project) LTE (Long Term Evolution) TDD system, as shown in Table 1, there are seven kinds of uplink-downlink configurations in total, where ‘D’ represents a downlink subframe, ‘U’ represents an uplink subframe and ‘S’ represents a special subframe. Since the special subframe includes a DwPTS, the special subframe can be used for downlink\nTABLE 1Uplink-downlink configurations in a TDD systemuplink-downlinksubframe serial number in aconfigurationrepeatradio frameserial numberperiod012345678905 msDSUUUDSUUU15 msDSUUDDSUUD25 msDSUDDDSUDD310 ms DSUUUDDDDD410 ms DSUUDDDDDD510 ms DSUDDDDDDD65 msDSUUUDSUUD\nIn a LTE system, which uplink-downlink configuration to be used is configured semi-statically by a base station, so it will occur that a configured uplink-downlink configuration does not match an instantaneous service type, and thus resources cannot be used efficiently, which is especially serious in a cell with a relatively small number of users. In order to solve this problem, a technology of applying a dynamic TDD subframe may be introduced to the LTE R11 Version, namely that several dynamic subframes, which are able to not only be used as uplink subframes but also be used as downlink subframes, are configured in a radio frame. Table 2 illustrates a schematic diagram of applying a dynamic TDD subframe in a radio frame, that is, the last two subframes of each half radio frame are dynamic subframes, wherein ‘F’ represents a dynamic subframe. In addition, subframes with subframe numbers 0, 1, 5 and 6 in a radio frame are all used for downlink, namely downlink subframes, and these four subframes may also be referred to as fixed downlink subframes. Subframes with subframe numbers 2 and 7 in a radio frame are both uplink subframes, and these two subframes may also be referred to as fixed uplink subframes.\nTABLE 2Subframe number in a radio frame0123456789DSUFFDSUFF\nHowever, in a scenario of applying the dynamic TDD subframe, the related art does not solve a problem how to determine whether a dynamic subframe is an uplink subframe or a downlink subframe, thus an information transmission in the scenario of applying the dynamic TDD subframe can not be implemented."}
{"text": "Many beneficial human manufacturing and energy production activities involve the transport of a heated fluid in a fluid conduit situated in a cold environment. Where the fluid is susceptible to solidification or becoming unmanageably viscous because of heat loss to the cold environment, prudent engineering practices include insulating pipes against passive heat loss to the environment and/or actively heating the fluid conduit along its length.\nHeat loss and its attendant consequences may become particularly severe where flow of the heated fluid through the conduit is interrupted. The conduit, containing the heated liquid gradually cools via heat loss to the environment and the fluid may solidify or become unmanageably viscous within the conduit. When flow is resumed, the thermally depleted matter within the fluid conduit may prevent or delay the resumption of fluid flow within the fluid conduit. Problems can be particularly severe when the heated fluid readily crystallizes on cooling, as is the case with relatively pure phenol (melting point 43° C., CAS No. 108-95-2) or otherwise forms solids on cooling. (See natural gas hydrates for example.)\nThus, heat retention within fluid conduits may be critical to the efficient operation of facilities in which a hot fluid susceptible to deleterious phase changes in response to heat loss to a cold environment is transported. There is at present a particular need for improved thermal control in subsea hydrocarbon production operations in which hot production fluids may undergo deleterious phase changes as a result of heat loss to the cold subsea environment. The present invention provides one or more embodiments enabling improved thermal control in such environments."}
{"text": "The present invention relates to an extreme ultraviolet (“EUV”) light spectrum measuring apparatus and an EUV light intensity calculating method used to evaluate an EUV light source (having a wavelength between 10 and 15 nm) for use with a projection exposure apparatus in a photolithography in manufacturing semiconductor devices.\nReduction projection exposures using ultraviolet have been conventionally employed to manufacture such fine semiconductor devices as a semiconductor memory and a logic circuit in the photolithography method. The critical dimension to be transferred by the reduction projection exposure is proportionate to a wavelength of light used for transfer, and inversely proportionate to the numerical aperture (“NA”) of a projection optical system. In order to transfer a finer circuit pattern, a shorter wavelength of used ultraviolet (“UV”) light has been promoted from an ultra-high pressure mercury lamp i-line with a wavelength of about 365 nm to a KrF excimer laser with a wavelength of about 248 nm and an ArF excimer laser with a wavelength of about 193 nm.\nHowever, the lithography using the UV light has the limits to satisfy the rapidly promoting fine processing to semiconductor devices, and a reduction projection exposure apparatus using the EUV light with a wavelength of about 10 to 15 nm much shorter than that of the ultraviolet has been developed to efficiently transfer a very fine circuit pattern of 0.1 μm or less.\nA development of the EUV light source for the reduction projection exposure apparatus is promoted parallel to a development of the reduction projection exposure apparatus. One illustrative EUV light source is, for example, a laser excited EUV light source as disclosed in Japanese Patent Application, Publication No. 2002-174700 (corresponding to U.S. Pat. No. 6,324,256). The laser excited EUV light source irradiates a highly intensified pulse laser beam to a target material put in a vacuum chamber, and generates the high-temperature plasma. Among the lights having various wavelengths emitted from thus generated plasma, the EUV light utilizes the light having a wavelength, for example, of about 13 nm. The target material uses a metallic thin film, an inert gas, a droplet, etc., and is supplied to a vacuum chamber by such a means as a gas jet. The EUV light emitted from the plasma is generally condensed into a condensing point by a rotationally elliptic condenser mirror, is introduced into a projection exposure apparatus after diverging from the condensing point, and illuminates a mask uniformly via an illumination optical system of the projection exposure apparatus. Alternatively, the EUV light may be introduced, as light collimated by the rotationally elliptical condenser mirror, into the projection exposure apparatus.\nIn the exposure apparatus using the EUV light, an optical element in the optical system that introduces the EUV light mainly uses an oblique-incidence total reflection mirror and a Si/Mo multilayer coating mirror as a mirror having an incident angle close to normal incidence. Since the multilayer coating mirror for the normal incidence has a high reflectance to the EUV light having a wavelength around 13.5 nm, the EUV light for projection exposure generally uses a wave range between 13.365 nm and 13.635 nm around the wavelength of 13.5 nm among the lights emitted from the EUV light source. The throughput of the projection exposure apparatus depends upon the absolute value of the EUV light intensity in the wave range between 13.365 nm and 13.635 nm, and an exposure apparatus has a higher productivity as the absolute value increases.\nThe EUV light source needs to emit the EUV light at a high intensity in the wave range between 13.365 nm and 13.635 nm among the lights having various wavelengths. The light intensity in the non-exposure wave range should be low, because this light turns to the heat after absorbed in the optical element in the optical system, and deteriorates the of the optical system.\nFor these reasons, it is important to previously recognize a spectrum of the light emitted from the EUV light source for the exposure apparatus.\nOne illustrative, conventional EUV light spectrum measuring apparatus is disclosed, for example, in Japanese Patent Application, Publication No. 10-073698. According to this reference, a diffraction grating forms a spectrum of part of the EUV light (or referred to as “soft-X-ray” in this reference) emitted from the plasma as a light emitting point, making the light intensity of each wavelength measurable.\nHowever, the above prior art can measure the spectrum of the light only in a specific divergent direction and cannot structurally measure the spectra in all the divergent directions. When the spectra differ according to divergent directions, the prior art cannot measure the spatial spectrum distribution or the intensity in the predetermined wave range, such as between 13.365 nm and 13.635 nm, in the entire divergent directions. In addition, the prior art disadvantageously cannot measure a total intensity in the predetermined wave range emitted from the light source."}
{"text": "1) Field of the Invention\nThis invention relates generally to fabrication of semiconductor devices and more particularly to the fabrication of an oxynitride layer for use as an anti-reflective coating (ARC) or as a stop layer having low leakage current.\n2) Description of the Prior Art\nThe semiconductor industry's continuing drive toward integrated circuits with ever decreasing geometries, coupled with its pervasive use of highly reflective materials, such as polysilicon, aluminum and metal silicides, has led to increased photolithographic patterning problems. Unwanted reflections from these underlying materials during the photoresist patterning process cause the resulting photoresist pattern to be distorted. The use of an anti-reflective coating (ARC) has become popular in the semiconductor industry to prevent distortion of photoresist patterns. An Anti-Reflection Coating (ARC) is a layer of electiomagnetic energy absorbing material formed subjacent to a photoresist layer to minimize pattern distortion caused by reflection. Oxynitride can be used to form anti-reflective coatings.\nImportantly, the inventors have found that to increase manufacturing efficiency, it would be advantageous in the semiconductor industry to use an ARC as a dielectric in subsequently defined semiconductor devices. However, the inventor's have found that oxynitride ARCs exhibit leakage currents that prevent their use as a dielectric in a semiconductor device.\nThe importance of overcoming the various deficiencies noted above is evidenced by the extensive technological development directed to the subject, as documented by the relevant patent and technical literature. The closest and apparently more relevant technical developments in the patent literature can be gleaned by considering the following patents.\nU.S. Pat. No. 5,639,687 (Roman et al.) discloses a process of making a Si-rich silicon oxynitride layer.\nU.S. Pat. No. 5,365,104 (Godinho et al.) shows a SiON process for making a silicon oxynitride passivation layer.\nU.S. Pat. No. 4,717,631 (Kaganowicz et al.) shows a method for making SiON having a specific low Hydrogen % range."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of signal processing. More specifically, the present invention relates to the processing of measured signals, containing a primary signal portion and a secondary signal portion, for the removal or derivation of either the primary or secondary signal portion when little is known about either of these components. The present invention is especially useful for physiological monitoring systems including blood oxygen saturation systems and pulserate measurement systems. The present invention further relates to a method and apparatus for signal processing of signals in order to compute an estimate for pulserate.\n2. Description of the Related Art\nSignal processors are typically employed to remove or derive either the primary or secondary signal portion from a composite measured signal including a primary signal portion and a secondary signal portion. For example, a composite signal may contain a primary signal portion comprising desirable data and a secondary signal portion comprising noise. If the secondary signal portion occupies a different frequency spectrum than the primary signal portion, then conventional filtering techniques such as low pass, band pass, and high pass filtering are available to remove or derive either the primary or the secondary signal portion from the total signal. Fixed single or multiple notch filters could also be employed if at least one of the primary and secondary signal portions exists at a fixed frequency band.\nIt is often the case that an overlap in frequency spectrum between the primary and secondary signal portions exists. Complicating matters further, the statistical properties of one or both of the primary and secondary signal portions may change with time. In such cases, conventional filtering techniques are ineffective in extracting either the primary or secondary signal. If, however, a description of either the primary or secondary signal portion can be derived, correlation canceling, such as adaptive noise canceling, can be employed to remove either the primary or secondary signal portion of the signal isolating the other portion. In other words, given sufficient information about one of the signal portions, that signal portion can be extracted.\nConventional correlation cancelers, such as adaptive noise cancelers, dynamically change their transfer function to adapt to and remove portions of a composite signal. However, correlation cancelers and adaptive noise cancelers require either a secondary reference or a primary reference which correlates to either the secondary signal portion only or the primary signal portion only. For instance, for a measured signal containing noise and desirable signal, the noise can be removed with a correlation canceler if a noise reference is available. This is often the case. Although the amplitudes of the reference signals are not necessarily the same as the amplitudes of the corresponding primary or secondary signal portions, the reference signals have a frequency spectrum which is similar to that of the primary or secondary signal portions.\nIn many cases, nothing or very little is known about the secondary and primary signal portions. One area where measured signals comprising a primary signal portion and a secondary signal portion about which no information can easily be determined is physiological monitoring. Physiological monitoring generally involves measured signals derived from a physiological system, such as the human body. Measurements which are typically taken with physiological monitoring systems include electrocardiographs, blood pressure, blood gas saturation (such as oxygen saturation), capnographs, other blood constituent monitoring, heart rate, respiration rate, electro-encephalograph (EEG) and depth of anesthesia, for example. Other types of measurements include those which measure the pressure and quantity of a substance within the body such as cardiac output, venous oxygen saturation, arterial oxygen saturation, bilirubin, total hemoglobin, breathalyzer testing, drug testing, cholesterol testing, glucose testing, and carbon dioxide testing, protein testing, carbon monoxide testing, and other in-vivo measurements, for example. Complications arising in these measurements are often due to motion of the patient, both external and internal (muscle movement, vessel movement, and probe movement, for example), during the measurement process.\nMany types of physiological measurements can be made by using the known properties of energy attenuation as a selected form of energy passes through a test medium such as a finger, shown schematically in FIG. 1.\nA blood gas monitor is one example of a physiological monitoring system which is based upon the measurement of energy attenuated by biological tissues or substances. Blood gas monitors transmit light into the test medium and measure the attenuation of the light as a function of time. The output signal of a blood gas monitor which is sensitive to the arterial blood flow contains a component having a waveform representative of the patient's arterial pulse. This type of signal, which contains a component related to the patient's pulse, is called a plethysmographic wave, and is shown in FIG. 2A as a curve s(t) 201. Plethysmographic waveforms are used in blood gas saturation measurements. As the heart beats, the amount of blood in the arteries increases and decreases, causing increases and decreases in energy attenuation, illustrated by a cyclic wave seen in the curve 201.\nTypically, a digit such as a finger, an ear lobe, or other portion of the body where blood flows close to the skin, is employed as the medium through which light energy is transmitted for blood gas attenuation measurements. The finger comprises skin, fat, bone, muscle, etc., as shown FIG. 1, each of which attenuates energy incident on the finger in a generally predictable and constant manner. However, when fleshy portions of the finger are compressed erratically, for example by motion of the finger, energy attenuation becomes erratic.\nAn example of a more realistic measured waveform is shown in FIG. 2B, as a curve M(t) 202. The curve 202 illustrates the effect of motion and noise n(t) added to the clean waveform s(t) shown in FIG. 201. The primary plethysmographic waveform portion of the signal M(t) is the waveform representative of the pulse, corresponding to the sawtooth-like pattern wave in curve 201. The large, secondary motion-induced excursions in signal amplitude obscure the primary plethysmographic signal s(t). Even small variations in amplitude make it difficult to distinguish the primary signal component s(t) in the presence of a secondary signal component n(t).\nA pulse oximeter is a type of blood gas monitor which non-invasively measures the arterial saturation of oxygen in the blood. The pumping of the heart forces freshly oxygenated blood into the arteries causing greater energy attenuation. As well understood in the art, the arterial saturation of oxygenated blood may be determined from the depth of the valleys relative to the peaks of two plethysmographic waveforms measured at separate wavelengths. Patient movement introduces motion artifacts to the composite signal as illustrated in the plethysmographic waveform illustrated in FIG. 2B. These motion artifacts distort the measured signal."}
{"text": "Present apparatus and methods for changing out carbon in an industrial carbon filter system, such as in a chemical process plant, include 50 to 2000 pound carbon canisters having openings at the top for receiving fresh carbon. Releasable attachment parts affixed to the canisters, such as nipples and flanges, are provided for connecting the canister into a process stream line for filtering the stream. A stream inlet nipple is usually provided near the top of the canister and a stream outlet nipple near the bottom.\nThe plant typically provides a roll-off box or the like for collecting spent carbon and sending it off to be revitalized and then re-used. The plant also usually provides a fork lift for moving canisters around the site. Typically a service provider is retained by the plant, however, to change out spent carbon and to refill the canisters with fresh carbon, upon demand and need. Spent carbon is shipped off-site to be revitalized. This spent carbon may be wet or dry.\nCurrent service provider systems and methods typically include a truck, manned by at least two people, the truck having a hopper and vacuum apparatus. At least two people are required to set up the hopper and run the vacuum operation. In operation, in the prior art, canisters with spent carbon filter material are disconnected from their process streams and brought to a central site, usually proximate the location of a roll-off box or the like at a plant facility. A service provider vacuum truck is driven to the site. The at least two service people set up the hopper as the well as the vacuuming equipment. The upper lid of the canisters are opened, each in turn, and spent carbon is vacuumed out from the canister and into the provider's hopper. When the hopper is filled, (a hopper typically holds about 1500 pounds of spent carbon,) vacuuming is stopped and the hopper is frequently moved by the fork lift to a roll-off box, dumped into the roll-off box and returned to the proximity of the truck for further vacuuming. In a less favored system, because it entails higher risks, the hopper dumps its contents into one or more spent carbon container bags at the canister site. The bags are then tied off and moved to a location for collecting and shipping off site. (The possibility of the deterioration of the bags and resulting contamination of the site prior to shipping causes this option to be slightly less favored.)\nThe instant invention comprises an improvement to these known prior art systems and methods. The instant invention, at least in certain key applications, has been tested to determine that it is more cost effective and requires less equipment and fewer people and takes less time.\nIn preferred key embodiments of the instant invention, a canister is provided with two valved orifices. While one orifice may always be provided on top of the canister for re-filling the canister with fresh carbon, preferably a second valved orifice is provided by the instant invention on lower parts of the canister or on the bottom. This second orifice is preferably situated on the bottom and preferably includes a ball valve or the like for controlling an approximately four to six inch diameter opening. The same roll-off box and forklift provided by the plant, as discussed above, can be used in the instant invention. As will be shown, however, the instant invention can also be utilized with existing canisters.\nIn key preferred embodiments of the instant invention, a truck arrives with one person to the site of the canisters. (The truck may include vibration equipment for difficult canisters.) In this preferred embodiment one person lifts, with the fork lift, each canister up and over the top of the plant's roll-off box or the like. (Alternate embodiments are also disclosed below.) The service person opens a lower orifice on the bottom of the canister using a ball valve, and the spent carbon drops by gravity into the roll-off box. Alternately, the service person might pivot the canister, to invert it over the roll-off box, such that the spent carbon dumps from the top orifice.\nThe instant inventor's tests indicate that while it takes two people 8-9 minutes to vacuum a thousand pounds of spent carbon from a canister, it only takes one person two minutes or less to dump a thousand pounds of spent carbon into a roll-off box. (Note: the canister could also be dumped into saturated carbon containers or bags instead of a roll-off box. Again, use of such bags is slightly disfavored, as mentioned above, for incurring higher potential contamination hazards.)\nThe instant invention eliminates the need for vacuuming equipment, for the second or more service person and for the intermediate step of filling and emptying the service provider's hopper.\nIn the odd situation when spent carbon might resist dropping out of a canister by gravity, tests have indicated that the techniques of shaking, scraping and/or vibration can be successfully employed. The inventor's tests have demonstrated that a forklift can affectively shake a canister to cause sections of spent wet carbon that failed to fall by gravity, to break off and fall out of a canister. The inside walls of a canister can also be scraped, or squeegeed, by the service person. In addition, a vibrator, such as a vibrator used with a cement mixer, can be connected or inserted in the canister, possibly with a hand held prong inserted inside of the canister, to help ensure that all spent carbon is dislodged.\nAgain, preferred novel canisters of key embodiments of the instant invention have a second bottom valved orifice, which opens and closes. Although these canisters should be more expensive than the canisters of the prior art, that cost should be more than recovered by the enhanced speed and lessened cost of dumping rather than vacuuming, and also by the reduced cost of manpower and equipment."}
{"text": "This invention relates to spread-spectrum communications and more particularly to a satellite broadcasting or paging system which communicates over the entire bandwidth of an existing frequency division multiple access (FDMA), or time division multiple access (TDMA), satellite system such as the proposed low earth orbiting (LEO) satellite communication system."}
{"text": "Refrigeration devices, such as, for example, refrigerators, freezers or wine storage cabinets, tend to become increasingly large to provide larger storage capacity. Regardless of their use, such refrigeration appliances have a body having an interior for receiving the objects or products to be cooled, and a door attached to the body. Refrigeration appliances always have an elastically deformable seal disposed between the body and the door to prevent warm air from entering the interior, and to prevent the cold present therein from escaping to the outside when the refrigeration appliance is closed. Consequently, the seal is of particular importance because it is also decisive for the power consumption of the refrigeration appliance. Due to the temperature difference between the environment and the interior of the refrigeration appliance, the warm air that has entered the interior of the refrigeration appliance after opening the door is cooled once the door is closed, causing a decrease in the volume of the air, as a result of which a vacuum develops in the interior of the refrigeration appliance and presses the door against the body with a force increased by the ambient air pressure. As a result of this naturally increased door-closing force, the door may be impossible to open or may be openable only with substantial force, especially with increasing size of the refrigeration appliances.\nIn order to reduce this force, German Utility Model DE 20 2000 011 427 U1 describes a refrigeration appliance having a body and a door closing the same, where an opening mechanism referred to as “drive means” acts on the door and is actuatable by a trigger sensor. According to the disclosure of this document, the trigger sensor may operate based on different principles, of which mechanical, acoustic or optical operating principles are mentioned by way of example. This document also mentions a force sensor which senses an actuation force acting on the door. Thus, the opening mechanism can be activated, for example, by slightly pushing against the door or slightly pulling on the door handle. The opening mechanism used here is an electric motor having bevel gearing.\nFurther, DE 20 2005 011 427 U1 also presents an opening mechanism which is composed of a driving gear wheel and a toothed rack and intended solely to assist the opening movement.\nIn German Patent Application DE 10 2006 061 083 A1, a sensor is described which is preferably disposed in the interior of the refrigeration appliance and, in accordance with the disclosure of this document, may take the form of both a tension sensor and a pressure sensor. This sensor, too, serves to sense movement of the door through pressure measurement, the signal being usable for activating an opening mechanism via a control circuit.\nThe approaches known heretofore have been found to have the disadvantage that the individual components have to be mounted at specific positions of the refrigeration appliance, which overall makes them complex to install. Some designs provide that the sensor is installed in the door, as a result of which electrical lines need to be run from the body into the door, which entails additional complexity. In sum, until now, there are no designs of opening mechanisms in refrigeration appliances that would be easy to install and, in addition, also satisfy aesthetic considerations."}
{"text": "It is known that effects of memory access latencies in a data processing system may be mitigated by moving elements of vector data into a local, high-speed memory known as a cache. The elements are moved, or prefetched, into the cache before they are needed so that they are readily available when requested. If the elements are in a predictable order, there is no theoretical limit to how far in advance the elements may be fetched. However, since the cache has a limited size, if elements are fetched too far in advance, prior elements may be displaced from the cache before they have been used. This can lead to the phenomenon of “thrashing”, where an element is prefetched and displaced multiple times before they are used. As a consequence, the performance of the data processing system may be worse than if no cache is used.\nOne approach to prevent thrashing is to prefetch elements directly from memory or indirectly via a cache into a sequentially ordered storage memory or queue. Once queued, elements remain in the queue until used, thus reducing thrashing. A disadvantage of this approach is that repeated accesses to an element either within or between vectors will result in the element being duplicated in the queue as well as in the cache. Additionally, if the element is evicted from the cache before subsequent prefetches, then the element must be fetched from memory again. In the worse case, this would result in a performance similar to having no cache.\nAnother approach is to provide prefetch instructions and place the burden on the programmer to use them in a manner that avoids thrashing. A disadvantage of this approach is that changes in the sparseness of the vector elements, the number of vectors, the memory access latency and even cache line replacement policy can require changes in the placement of the prefetch instructions."}
{"text": "In typical heavy-duty vehicle applications, including those with hybrid drive systems, a rotating internal combustion engine includes multiple gear and/or pulley and belt power take-offs (PTOs) that operate the vehicle subsystems and accessories. As a result, turning off the engine causes the vehicle subsystems and accessories to be turned off.\nIt is desirable to eliminate engine idling at vehicle stops to, among other things, increase fuel economy, minimize noise, and minimize engine exhaust emissions pollution to improve the quality of the operating environment. This is especially true for transportation and delivery vehicles such as, but not limited to, urban transit buses, refuse collection trucks, and local package freight pick up and delivery vans that may experience hundreds of stops during daily operation.\nA driver could manually turn off and turn on an engine when stopped and restarted; however, in addition to the problem of the vehicle subsystems and accessories not being operation, a typical electric starter motor for the internal combustion engine would wear out rather quickly because it is not designed for the hundreds of stop-starts per day of transportation, collection, and delivery vehicles."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic garnet single crystal and a method for producing the same as well as an optical element using the same.\n2. Description of the Related Art\nFaraday rotator is an optical element with a function to rotate the polarization plane of transmitting light and is used in optical devices such as optical isolator, optical attenuator, optical circulator and photo-magnetic field sensor for communication. Faraday rotator is generally prepared by using a plane-like bismuth (Bi)-substituted rare-earth iron garnet single crystal. The Bi-substituted rare-earth iron garnet single crystal is grown by the liquid phase epitaxial (LPE) process as one of flux processes. Single crystals are grown at atmospheric pressure by the flux processes.\nIn growing the Bi-substituted rare-earth iron garnet single crystal by the solution process (the LPE process) such as flux processes, generally, lead oxide (PbO), bismuth oxide (Bi2O3) and boron oxide (B2O3) are used as the solvents so as to stably grow such garnet single crystal while the solvents are maintained at their supersaturation states. During the growth of the magnetic garnet single crystal, therefore, a small amount of lead (Pb) contaminates in the resulting crystal. In Faraday rotators for use in optical devices for communication, magnetic garnet single crystals at a Pb content “y” of about 0.03 to 0.06 in the chemical formula Bi3-x-yM11xPbyFe5-z-wM12zM13wO12 have been used conventionally. See Patent Reference 1: JP-A-2001-044026; Patent Reference 2: JP-A-2001-044027; and Patent Reference 3: JP-B-6-046604.\nFollowing the upsurge in the recent environmental protection movement, however, efforts are now directed toward the reduction of the content of Pb as an environmental load substance in any of industrial products. Therefore, a trace amount of contaminating Pb in magnetic garnet single crystals grown by the LPE process draws concerns as a factor of environmental pollution. Therefore, it is necessary to reduce or eliminate the amount of Pb contained in magnetic garnet single crystals as materials constituting Faraday rotators."}
{"text": "The washing of preserved rawhides and fur skins before tanning has the purpose of removing dirt, blood, dung, preservatives, fat and water-soluble proteinacous compounds. This is done mostly in a paddle tub or in the tanning tumbler. Frequently service water is used which has a hardness of about 15.degree. dH (degrees german hardness). The washing and cleaning solutions employ the following aids:\n(a) Surface-active substances of an anionic or nonionic nature, which not only clean the skins, but also improve at the same time the feel of the fur.\n(b) Fat solvents in emulsified form, such as hydroaromates (hydrogenated aromatic hydrocarbons) or petroleum-hydrocarbons.\n(c) Inorganic salts, like common salt, which help to improve the washing effect of the tensides by an electrolytic effect. A high pH of 8.5 or over is undesired and can lead to damage of the skins and rawhides. The use of the inorganic salts leads to a high salt content of the waste waters, which is subject to increasing criticism.\nA further problem in washing rawhides and skins, particularly when treating raw materials with a high fat content, is that the cleaned product can become fatted again (refatted) and the apparatus can be contaminated in the dilution stage on rinsing. In other words the suspended fat tends to deemulsify and deposit on dilution during rinsing. An improved stabilization of the wash liquors is thus desirable. Furthermore questions of sewage treatment and the resulting costs are becoming important and will become more and more important in the future."}
{"text": "Disadvantages mentioned above are eliminated significantly by the method of authenticity and/or agreement approval at the direct debits through the separate identifier according to this invention, of which bedrock is based on the fact that after activating the payment process or payment process preparation in the mobile communication device, the device sends alphanumerical chain into the identifier approached to the mobile communication device. The identifier is approached in the distance shorter than 10 cm, mostly it may touch the mobile communication device directly. However, the contact or the contact point setting is not essential or necessary since the communication between the mobile communication device and the identifier is wireless. At this point, the identifier is by means of energy supplied contact free from the electromagnetic field of the mobile communication device, receives alphanumerical chain and realizes its processing in the form of an electronic signature. The received alphanumerical chain is signed electronically in the identifier and such established signed alphanumerical chain is sent back to the mobile communication device where the correctness is checked. If the correctness is verified and at the same time the owner approves the payment transaction, the mobile communication device will be approached to the distance, suitably less than 10 cm from the reading unit of the payment terminal.\nIt is advantageous if the electronic signature of the received alphanumerical chain is made by the process with the help of a private key saved in the memory of the identifier. To the energy supply of the identifier circuits is used direct and/or while approaching, accumulated energy of the electromagnetic field of the mobile communication device, advantageously electromagnetic field created by transmitting unit of the mobile communication device which aim is to communicate with the identifier.\nThe invention enables to use cryptography of the authentication preserving the energy passivity of the separate identifier. The principal advantage is the high level of security at the satisfactory user comfort. Entering PIN code mostly consisting of four digits can be according to this invention, replaced by apposing the identifier working with more-bit-chain which is variable since the approving code transmitted from the identifier into the mobile communication device is changed according to the cryptography taking place in the identifier processor at every authentication process. By this, the higher level of security is achieved than at the RFID identifiers and the advantage of its passivity is remained. Such energy passivity enables to reduce the size of the identifier since there is no necessity to use own energy source by which the user comfort is increased as the user does not have to bother with identifier's charge or condition.\nIn a favorable configuration the mobile communication device may use at the payment processing and authentication a remote processing server to which it is connected via general mobile network, mainly the kind of GSM or GRPS net.\nThe bedrock of this invention is based on the system of authentication and/or agreement approval at the direct debits which concerns payment terminals, a separate identifier and a mobile communication device communicating with the payment terminal via contact free communication channel where the mobile communication device, preferably mobile phone, contains transmitting and receiving unit to allow the contact free communication with the identifier and where the identifier contains a processor for electronic signature of the received alphanumerical chain. Further more, the identifier contains transmitting and receiving unit to allow the communication with the mobile communication device, a block transforming the electromagnetic field into the electric energy and a memory. Components of the identifier are, from the energy point of view, supplied by electromagnetic field of the mobile communication device either directly by immediately gained energy or partially by energy accumulated during the approaching to the mobile communication device at the relevant payment process.\nIn a possible configuration the system includes a remote processing server connected to the mobile communication device through general mobile network, preferably via GSM or/and GPRS net. From the point of view of compatibleness with the existing, mass-spread devices and standards, it is suitable if the mobile communication device consists of a mobile phone, advantageously a mobile phone with NFC communication unit.\nDisadvantages mentioned in the Present Technology Status are eliminated significantly by the identity identifier and/or agreement approval at the direct debits through communicating contact free with the mobile communication device, principally a mobile phone which is connected contact free to the payment terminal according to this invention, of which bedrock is based on the fact that it consists of a processor for electronic signature of the received alphanumerical chain, transmitting and receiving unit to communicate with the mobile communication device. Communication is principally based on receiving the alphanumerical chain and transmitting electronically signed alphanumerical chain. The identifier further more contains a memory, a block transforming the electromagnetic field into the electric energy. Transmitting and receiving unit and a block transforming the electromagnetic field into the electric energy are connected to the processor. Processor is also connected to the memory. Basically, all the units of the identifier are, from the energy point of view, supplied by electromagnetic field of the mobile communication device.\nIn an advantageous configuration, the identifier of identity and/or approval contains NFC chip and the memory contains a private key for electronic signature of alphanumerical chain received from the mobile communication device.\nUtility attributions increase a configuration, where a part of the memory is reserved for the personal data of the user. This part of memory is adjusted to store the personal user data separately from the private key. In such case the identifier may be used like a health insurance card, identity card and so on. Rightness of data demand is evaluated in the processor of the identifier.\nTo increase the user comfort, the identifier may be located in a pendant and/or a key ring and/or a label and/or a beading.\nThe invention enables to higher the level of security and comfort of authenticity and direct debit approval since the user does not have to remember his PIN code. At the same time the invention increases the process of direct debit transaction since the mobile communication device functions as a wallet without any delay caused by the PIN code entering."}
{"text": "1. Technical Field\nThe present invention relates to a method and apparatus for fabricating an optical identification element; and more particularly to a method and apparatus for fabricating a holographic optical identification element using a lithographic technique, as well as the holographic optical identification element itself."}
{"text": "This invention relates to a process for the catalytic hydrogenation of plant-based steroid compounds. More particularly, this invention relates to the metal-catalyzed hydrogenation of phytosterols and to the recovery and regeneration of the metal catalyst species used in such metal-catalyzed hydrogenation.\nSteroids are macromolecular organic compounds that, structurally, comprise a cyclopentanephenanthrene skeleton. Many of these compounds are alcohols and are identified as sterols. When combined to form esters, these compounds are called xe2x80x9cceridesxe2x80x9d.\nCholesterol, a sterol, has been known for some time. It has been reputed as being involved in the formation of biliary calculi, a theory which has been confirmed by proof of its involvement in circulatory disorders, and more particularly in the hardening of the arteries.\nNonesterified cholesterol is the main constituent of fatty substances having an animal origin. These animal fats are present in most of our foodstuffs and constitute an important source of cholesterol. When ingested in excess, they may breakdown, releasing cholesterol and the fatty acids which, in turn, may cause serious cardiovascular diseases.\nHowever, fatty substances having a plant origin do not contain cholesterol, but instead comprise a mixture of cholesterol-like materials called phytosterols, i.e., they are unsaturated, plant-based compounds comprising the cyclopentanephenanthrene skeleton. The best known of such compounds include: stigmasterol, sitosterol and ergosterol.\nThe structure of the molecule on which the animal and plant steroids is based is shown below: \nThe C 3 carbon atom is hydroxyl-substituted. The carbon atom at position 10 bears a methyl group.\nThe stereochemistry of the substituents at C 3 and C 10 is important. The steroid ring is a planar structure and, for cholesterol, the two groups at C 3 and C 10 lie above the plane of the ring, i.e., they are in the xcex2 conformation. The hydrogen atom at position C 6 (see the double bond between C 5 and C 6) lies below the plane, i.e., it is in the a position. Cholesterol\"\"s phytosterol look-alike, sitosterol, is substantially identical to cholesterol but, in addition bears an ethyl group at C 24.\nIn the reduction of these steroid-based structures, it is important to retain the stereochemical conformation of the starting material. Such reductions have typically been carried out with varying degrees of success, using hydrogen in the presence of supported metallic catalysts. See, for example, U.S. Pat. No. 3,865,853. Such catalysts include iron, cobalt, nickel, palladium, platinum, copper, silver and gold.\nAs disclosed in U.S. Pat. No. 6,147,235, steroids have been hydrogenated in the presence of catalytic amounts of nickel black, Raney Nickel, and nickel metal embedded in inorganic supports. However, these catalysts have not proven sufficiently selective for hydrogenation of the steroid double bond.\nThe \"\"235 patent also shows that steroids have been hydrogenated using noble metal catalysts, principally platinum and palladium. According to this patent, these catalysts have been supported on inorganic substrates, on carbon or used as metal blacks. The best conversions have been obtained when such noble metal catalysts are bound to activated carbon. The noble metal of choice is palladium and one of the most successful of the catalyst species has been shown to be palladium on carbon.\nHowever, prior art catalysts used for the reduction of steroids have proven to be disadvantageous because of difficulties of separation of the catalyst powder from the reaction mass after the hydrogenation process, the inflammability of the catalyst and the catalyst is not recyclable.\nAs shown on U.S. Pat. No. 6,147,235 palladium bound to an organic substrate (a polyolefin or halo polyolefin) also has been used successfully. The process for manufacturing such a catalyst is difficult and the expense of such a catalyst species has limited its use.\nIt is the object of the present invention to provide a process for catalytically hydrogenating unsaturated phytosterol compounds. The catalyst used for such hydrogenation is an alumina-supported transition metal selected from the group consisting of nickel, palladium, platinum, ruthenium, rhodium and mixtures thereof.\nThe process of the present invention is as follows.\nAn admixture of a solvent, an alumina-supported transition metal catalyst and a material comprising at least one of the unsaturated phytosterol compounds is first formed. This admixture is hydrogenated for a time and at a temperature and pressure sufficient to cause the hydrogenation of at least one of said unsaturated phytosterol compounds.\nThe alumina-supported transition metal catalyst is next separated from the reaction (the hydrogenated) mixture that was first formed. This separated alumina-supported transition metal catalyst is then admixed with a new mixture containing a solvent and a material comprising at least one unsaturated phytosterol compound. The resulting reaction mixture is hydrogenated for a time and at a temperature and pressure sufficient to cause the hydrogenation of at least one of said unsaturated phytosterol compound. The alumina-supported transition metal catalyst is then separated from the hydrogenated mixture. It may be used again for subsequent hydrogenations of unsaturated phytosterol compounds.\nThe hydrogenation process disclosed herein can also be applied to the hydrogenation of fatty acids derived from phytosterols, e.g., tall oil, vegetable oil and their distillation products."}
{"text": "Garbage, waste, or any other materials which are compacted for transportation require the use of fixed walled containers or trailers which allow the transportation of materials in compacted form or compacting of the material in the container or trailer itself. These types of trailers cannot backhaul any bulk goods on the way back from the unloading site, as their fixed and closed top does not allow loading from the top using a conventional loader, similar to an open-top trailer.\nWith the rising cost of fuel, the disposal of garbage becomes increasingly expensive, especially when the trailer cannot backhaul any goods to reduce the cost of waste disposal, and thus, returns empty from the unloading site.\nOn the other hand, open top trailers, even those which include a cover, cannot be used for compacting or for the transport of loosely compacted material.\nUnited States Patent Application Publication No. 2007/0126209 (Wolfe) describes a trailer that converts into a shed. The trailer of Wolfe includes a convertible roof comprising two panels. However, the panels cannot be locked, and the walls are not rigid enough to withstand compressed materials.\nEuropean Patent Application No. 1,120,363 and United Kingdom Patent Application No. 2,429,006 (Sumpter) describe a trailer for the collection of waste comprising a chassis and a container. The container comprises a roof including an aperture closeable by a cover that slides over the aperture by means of a channel assembly. The cover may be a partly or fully retractable flexible sheet.\nThe structure of Sumpter does not allow the trailer to be used for transporting loosely compressed materials because the materials are distributed in the container using reciprocating slats, not a compactor. Furthermore, the cover disclosed in Sumpter is flexible and would thus not be able to withstand any pressure from within the container. Sumpter clearly does neither teach nor suggest using the container disclosed for transport of loosely compacted materials.\nNeither Wolfe nor Sumpter deal with loosely compacted materials, and thus, fail to recognize the structural challenges created by internal pressures in a trailer upon compression of loose materials therein.\nIt is, therefore, desirable to provide a trailer that can load and contain loosely compacted material, and which can also backhaul bulk goods on the way back from the unloading site in order to reduce the cost of transportation."}
{"text": "Conventionally, an analyzer for suctioning a reagent from a plurality of reagent containers mounted on a reagent table and analyzing a sample is known. The user needed to remove the reagent container to be replaced from the reagent table and mount a new reagent container to a predetermined position of the reagent table in such analyzer when reagent replacement was necessary. However, the user had to put his/her hand inside the accommodating unit of the reagent container in such analyzer, and thus the replacement task of the reagent was troublesome.\nAn analyzer that facilitates reagent replacement by the user is thus known (see e.g., Japanese Laid-Open Patent Publication No. 4-36658 and Japanese Laid-Open Patent Publication No. 2005-37171).\nAn analyzer including a rotatable reagent storage in which reagent bottles are arranged, a reagent bottle supply mechanism for supplying the reagent bottle to the reagent storage, and a reagent bottle discharge mechanism for discharging the reagent bottle from the reagent storage is disclosed in Japanese Laid-Open Patent Publication No. 4-36658. Such analyzer enables automatic replacement of the reagent bottle during analysis.\nAn analyzer including a reagent disc for installing a reagent, a replenishing reagent cabinet arranged at a position different from the reagent disc, and a reagent holding means for conveying the reagent from the replenishing reagent cabinet to the reagent disc is disclosed in Japanese Laid-Open Patent Publication No. 2005-37171. Such analyzer facilitates reagent replacement.\nHowever, the analyzer disclosed in Japanese Laid-Open Patent Publication No. 4-36658 has a complicating configuration since a belt line mechanism for conveying the reagent from a reagent bottle set yard to the reagent storage, and a belt line mechanism for conveying the reagent from the reagent storage to a reagent bottle discharge yard are necessary. Furthermore, in the analyzer disclosed in Japanese Laid-Open Patent Publication No. 4-36658, a space for arranging the reagent bottle set yard, the belt line mechanisms, and the reagent bottle discharge yard is necessary, which enlarges the device.\nThe analyzer disclosed in Japanese Laid-Open Patent Publication No. 2005-37171 has a complicating configuration since a mechanism (reagent holding means) for conveying the reagent in three-dimensional directions (XYZ directions) is necessary to convey the reagent from the replenishing reagent cabinet to the reagent disc. Furthermore, in the analyzer disclosed in Japanese Laid-Open Patent Publication No. 2005-37171, a space for arranging the reagent holding means and the replenishment reagent cabinet is necessary, which enlarges the device."}
{"text": "During seismic surveying, a seismic source is repeatedly actuated and seismic receivers, such as hydrophones in marine seismic surveying, receive energy direct from the sources and reflected from various boundaries or interfaces. In the case of marine seismic surveying, energy propagates into the earth and is reflected back to the hydrophones from subterranean boundaries or interfaces, for instance between strata of different types. The recorded seismic data can be processed to reveal information about the structure of the earth in the area being surveyed. However, such reflections are contaminated by other reflection paths. For instance, energy from the sources is reflected at the sea surface directly to the hydrophones. Also, energy can be reflected more than once between the sources and the receivers. Such multiple reflections can take place within the earth. Also, energy initially traveling downwards from the sources can be reflected upwardly and then downwardly again from the surface of the sea before arriving at the hydrophones. Reflections of this type are referred to as \"free-surface multiple reflections\". Free-surface multiple reflections can be classified according to their order, which is equal to the number of reflections from the free-surface. Thus, first order free-surface reflections comprise energy initially traveling downwardly from the sources (as opposed to \"ghosting\" where energy travels upwardly and is reflected from the free-surface), is reflected upwardly from the sea bed or a boundary below the sea bed, and is then reflected downwardly from the free-surface to the hydrophones. Second order free-surface multiple reflections undergo two downward reflections from the sea-surface before being detected by the hydrophones, and so on.\nThe invention makes use of a general concept of multiple attenuation using the inverse scattering or scattering Born series. For this general concept reference is made to the published International patent application WO 95/10787 and the British patent application No. 9426255.7 (published as GB-A-2296567).\nFrom those documents it is known how to derive the scattering Born series for removing free surface multiples: EQU D.sub.P =D.sub.0 +AD.sub.1 +A.sup.2 D.sub.2 +A.sup.3 D.sub.3 + . . . ,[1]\nwherein D.sub.P denotes the signals without free-surface multiples and A denotes the inverse of the source signature. The first term of the scattering series, D.sub.0, is the actual seismic data, the second term, D.sub.1, removes first order free-surface multiples, the following term, D.sub.2, removes second order free-surface multiples, and so forth.\nIt is an object of the present invention to improve and utilize the above method for sea bottom seismic acquisition operations."}
{"text": "Adequate sexual function is a complex interaction of hormonal events and psychosocial relationships. There are four stages to sexual response as described in the International Journal of Gynecology & Obstetrics, 51(3):265-277 (1995). The first stage of sexual response is desire. The second stage of sexual response is arousal. Both physical and emotional stimulation may lead to breast and genital vasodilation and clitoral engorgement (vasocongestion). In the female, dilation and engorgement of the blood vessels in the labia and tissue surrounding the vagina produce the \"orgasmic platform,\" an area at the distal third of the vagina where blood becomes sequestered. Localized perivaginal swelling and vaginal lubrication make up the changes in this stage of sexual response. Subsequently, ballooning of the proximal portion of the vagina and elevation of the uterus occurs. In the male, vasodilation of the cavernosal arteries and closure of the venous channels that drain the penis produce an erection. The third stage of sexual response is orgasm, while the fourth stage if resolution. Interruption or absence of any of the stages of the sexual response cycle can result in sexual dysfunction. One study found that 35% of males and 42% of females reported some form of sexual dysfunction. Read et al, J. Public Health Med., 19(4): 387-391 (1997).\nIn both pre-menopausal and menopausal females, sexual dysfunction can include, for example, sexual pain disorders, sexual desire disorders, sexual arousal dysfunction, orgasmic dysfunction, dyspareunia, and vaginismus. Sexual dysfunction can be caused, for example, by pregnancy, menopause, cancer, pelvic surgery, chronic medical illness or medications.\nErectile dysfunction is a widespread disorder that is thought to affect about 10% to 15% percent of adult men. A number of causes of erectile insufficiency, in addition to anatomical deficiencies of the penis or scrotum that preclude an erection sufficient for vaginal penetration, have been identified. Causes of erectile dysfunction can be categorized as psychogenic, neurogenic, endocrinologic, drug-induced, or vasculogenic and in any individual suffering from erectile dysfunction there may be more than one cause.\nPsychogenic impotence is often the result of anxiety or depression, with no apparent somatic or organic impairment Neurogenic impotence may arise from, for example, surgery or a pelvic injury, involving the nervous system affecting the penis. Erectile dysfunction which is endocrinologic in origin is most often associated with the disorders hypo- or hypergonadotropic hypogonadism and hyperprolactinemia.\nVasculogenic impotence is thought to be the most frequent cause of impotence accounting for approximately fifty percent of all cases of organic impotence. In these cases, the erectile dysfunction may be attributed to alterations in the flow of blood to and from the penis. Atherosclerotic or traumatic arterial occlusive disease to the arteries which supply blood to the penis can lead to a decrease in the rigidity of the erect penis as well as increase the time to achieving maximal erection. In still other cases, there is leakage from veins in the penis such that sufficient pressure for an erection can be neither obtained nor maintained.\nThere is also a high incidence of erectile insufficiency among diabetics, particularly those with insulin-dependent diabetes mellitus. Erectile dysfunction in diabetics is often classified as \"diabetogenic,\" although the underlying dysfunction is usually neurogenic, but may be vasculogenic or neurogenic and vasculogenic. About half of diabetic males suffer from erectile insufficiency, and about half of the cases of neurogenic impotence are in diabetics.\nErectile insufficiency is sometimes a side effect of certain drugs, such as beta-antagonists that are administered to reduce blood pressure in persons suffering from hypertension, or drugs administered to treat depression or anxiety. Excessive alcohol consumption has also been linked to erectile insufficiency. These forms of erectile insufficiency may be regarded as a subset of neurogenic or psychogenic insufficiency.\nA number of methods to treat impotence are available. These treatments include pharmacological treatments, surgery and, in cases of psychogenic dysfunction, psychological counseling is sometimes effective. Psychogenic impotence often can be cured by counseling coupled with a demonstration to the patient that he is capable of having a full erection by inducing such an erection from one to a few times in the patients. Insufficiency due to excessive alcohol consumption is sometimes cured by reducing or eliminating such consumption.\nIn the rare cases where the insufficiency is untreatable because of venous leakage, surgery can usually be used to repair the venous lesion and thereby either cure the insufficiency or, if there remains an erectile insufficiency after repair of the venous lesion, render the insufficiency amenable to treatment by pharmacological methods. Also, penile implants, which provide a mechanic means to produce an erection sufficient for vaginal penetration, are widely used to treat impotence. In recent years, implants have been used, especially in cases where pharmacological intervention is ineffective, which are usually cases of severe vasculogenic impotence. Treatment of impotence with penile implants, however, entails serious disadvantages. Such treatment requires surgery and necessitates total destruction of the erectile tissues of the penis, forever precluding normal erection.\nPharmacological methods of treatment are also available. Such methods, however, have not proven to be highly satisfactory or without potentially severe side-effects. Papaverine is now widely used to treat impotence, although papaverine is ineffective in overcoming impotence due, at least in part, to severe atherosclerosis. Papaverine is effective in cases where the dysfunction is psychogenic or neurogenic and severe atherosclerosis is not involved. Injection of papaverine, a phosphodiesterase inhibitor and a smooth muscle relaxant, or phenoxybenzamine, a non-specific antagonist and hypotensive, into corpus cavemosum has been found to cause an erection sufficient for vaginal penetration, however, these treatments are not without the serious and often painful side effect of priapism. Also, in cases where severe atherosclerosis is not a cause of the dysfunction, intracavernosal injection of phentolamine, an .alpha.-adrenergic antagonist, causes an erection sufficient for vaginal penetration. The resulting erection is one of significantly shorter duration than that induced by intracavernosal injection of papaverine or phenoxybenzamine and often times is of such short duration that satisfactory sexual relations are difficult or impossible. As an alternative or, in some cases an adjunct to phosphodiesterase inhibition or .alpha.-adrenergic blockade for the treatment of erectile dysfunction, prostaglandin E1 (PGE1) has been administered via intracavernosal injection. A major side effect frequently associated with intracorprally delivered PC;E1 is penile pain and burning.\nThus, there is a need in the art for treatments of male and female sexual dysfunctions, including treatments without the undesirable side effects of those agents currently used. The present invention is directed to these, as well as other, important ends."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor laser devices using gallium nitride-based semiconductor and optical information reproduction apparatuses using the semiconductor laser device as a source of light, and particularly to an improvement in the device's end surface mirror.\n2. Description of the Related Art\nNitride based semiconductor materials represented by GaN, InN, AlN and mix crystal semiconductor thereof are used to build a prototype of a semiconductor laser device emitting light in a range from blue color through ultraviolet. FIG. 10 shows a nitride semiconductor laser device lasing with a wavelength of 405 nm, as reported by Masaru KURAMOTO et al., Japanese Journal of Applied Physics, vol. 38, pp. L184-L186, 1999. This semiconductor laser device includes a substrate 901 of n-GaN (of 100 μm in thickness) and thereon a lower clad layer 902 of n-Al0.07Ga0.93N (of 1 μm in thickness), a lower guide layer 903 of n-GaN (of 0.1 μm in thickness), a triple quantum well active layer 904 of In0.2Ga0.8N (of 3 nm in thickness) /In0.05Ga0.95N (of 5 nm in thickness), a cap layer 905 of p-Al0.19Ga0.81N (of 20 nm in thickness), an upper guide layer 906 of p-GaN (of 0.1 μm in thickness), an upper clad layer 907 of p-Al0.07Ga0.93N (of 0.5 μm in thickness) and a contact layer 908 of p-GaN (of 0.05 μm in thickness), successively stacked, with an electrode 909 provided thereon and thereunder. Furthermore, an end surface mirror is provided by cleavage. This laser device has active and guide layers posed between clad layers to have a waveguide structure. The active layer emits light, which is in turn confined in this waveguide structure, and the end surface mirror functions as laser resonator mirror to provide a lasing operation.\nThis conventional semiconductor laser device as described above, however, is disadvantageous, as follows: the present inventors fabricated a semiconductor laser device having the aforementioned structure and have found that in a waveguide portion at an end surface mirror an epitaxially grown layer (each semiconductor layer upper than substrate 901) can have a cross section, a region in a vicinity of a front surface (i.e., a surface opposite to the substrate) in particular, failing to have a satisfactory cleavage plane. Herein an “unsatisfactory cleavage plane” means that an epitaxially grown layer forming a waveguide has a step offset from a vertical cross section or there exists a plane different in angle from a cleavage plane of the GaN substrate or there exists a rolling portion. The present inventors observed this portion in detail and have found that in particular at the active layer, cracking varies, and from the active layer through the front surface, the cracking is often observed different than in the remainder.\nA semiconductor laser device has an end surface mirror provided by externally exerting a force along a scratch on a wafer to warp and thus divide the wafer at a predetermined position. As such, the wafer particularly has a surface and therearound receiving a large magnitude of force exerted thereto and thus considered as cracking in a more complicated manner than the interior of the wafer. Possibly for the aforementioned semiconductor based semiconductor laser device this issue was not noted in designing the structure of a layer provided in a vicinity of an active layer formed in a vicinity of the wafer surface and it is thus believed that cracking tends to vary at the active layer of InGaN.\nThe conventional technology can thus fail to provide a satisfactory end surface mirror. Thus not only does mirror reflectance variation vary a threshold value, differential efficiency and other similar device characteristics but also a light radiation surface roughened prevents a far field pattern (FFP) from being smoothly unimodal, and a split peek, ripple or other impaired optical characteristics would be introduced. Such a FFP abnormality is not preferable as it results in insufficiently condensed light in application for example to an optical pickup in particular and in an extreme case causes stray light."}
{"text": "In photography, an image is recorded onto film and the image on the film is then developed into a photographic reproduction of the image. Usually the film contains an emulsion on which the image is recorded. In the standard photodeveloping process, the exposed film is processed in a developer which contains primary developing agents which reduce silver in the film emulsion to produce an image thereon.\nIn high resolution microfilm, the film emulsion is loaded with a high level of accutance, antihalation and sensitization dyes.\nWhen film is developed in a film processor, oily dye residues are released from the film emulsions. These oily dye residues build up over time. If these oily dye residues do not remain in solution, they will stain microfilm that is subsequently developed in the film processor or build up on the surfaces of the film processor, plugging orifices or otherwise adversely affecting the performance of the film processor.\nOne solution to the dye residue buildup problem is to change the developer solution in the film processor frequently. This, however, escalates the cost of photodevelopment and is not the preferred alternative. Furthermore, closed loop film processors are designed so that they will not have to be shut down prior to having developed a certain amount of film. Shutting down a closed loop process during a development cycle reduces the efficiency and economic advantages of the closed loop process.\nPreferably, an emulsifier is used to solubilize the oily residues in the photodeveloper to prevent the oily dye residues from adhering to the film or the components of the film processor. Standard emulsifiers such as disodium mono- and didodecyl diphenyl oxide disulfide anionic emulsifier or sodium alpha olefin sulfonate foam are not suited for this purpose, especially when the photodeveloping solution is continuously circulated through the film processor. These standard emulsifiers foam excessively when subjected to the constant recirculation of the developing solution in this type of film processor. The very alkaline conditions of the photodeveloper solutions in this type of chemical environment also adversely affect the stability of many emulsifiers.\nThe foam generated by these standard emulsifiers in photodeveloping solutions is very undesirable. Conventional surfactants are therefore of limited usefulness for preventing the build up of oily dye residues in the chemical environment required by film development, especially in closed loop film processors.\nBrightening agents have also been incorporated into photodeveloping solutions in order to control, reduce or eliminate the problem of dye stains on developed film. U.S. Pat. No. 4,264,716 to Vincent et al. incorporates brightening agents into photodeveloping compositions to reduce staining.\nHowever, brightening agents merely cover up or correct dye stains on photographs. Brightening agents do not solubilize the oily dye residues in the developer solution to prevent them from staining the film or adhering to components of the film processing tanks.\nTypical brightening agents include stilbene type brightening agents. Stilbene-containing compositions have been used to reduce or eliminate stains on film caused by the build up of oily dye residues in the developer solution.\nA published disclosure, Cappel, C.R. and Purol, M.D., \"Method for Reducing Spectral Sensitizing Dye Stain in Photographic Elements,\" Research Disclosures, No. 20733 (Kenneth Mason Publications, Ltd, Jul. 10, 1981) describes adding a water soluble stilbene, bis-4, 4'-s-triazyl amino stilbene sulfonic acid or mixtures of water-soluble stilbenes and nonionic surfactants, into a developing solution. The build up of dibenzooxazole carbocyanine sensitizing dyes in the gel matrix or within a coupler dispersion in the photographic emulsion layers was said to be avoided by the addition of the aqueous amino stilbene sulfonic acid solution. Dibenzooxazole carbocyanine sensitizing dyes extend the sensitivity of photographic emulsions to electromagnetic radiation which is beyond the wave lengths of visible light.\nU.S. Pat. No. 4,587,195 to Ishikawa et al. discloses a method of processing a silver halide photographic light-sensitive film material. The method is said to be economical because the amount of photodeveloping solution used in the process is reduced.\nIn Ishikawa et al. the reduction in the amount of photodeveloping solution used is accomplished by not changing the photodeveloping solution in the film processor as often. Since the photodeveloping solution is not changed as often, the problems associated with the accumulation of oily dye residues on film processing equipment become more acute.\nIn Ishikawa et al., a combination of two different triazyl stilbene-type brightening agents must be added to the color developer to avoid the adverse effects caused by the build up of oily dye residues therein. The brightening agents are present in the photodeveloper in an amount of about 0.3 gram to about 10 grams per liter of color developer solution.\nIf the maximum fluorescent wavelengths of the fluorescent spectra of the two triazyl stilbenes used in Ishikawa et al. are less than 4 m.mu. apart, the triazyl stilbene solution will not remove the dye stains, however. If only one of the two types of stilbene brightening agents is used or if the difference in the maximum fluorescent wavelength of the two triazyl stilbenes is not at least 4 m.mu. , build up of oily dye residues in the developer solution remains a problem. If the amount of triazyl stilbene in solution exceeds 10 grams per liter, the film's development will be adversely effected and fluorescent quenching or an unfavorable photographic performance such as non-color may result.\nFrom the foregoing it is apparent that the type of stilbene used to reduce dye stains on films is linked to the type of dye residue in the process. The type of dye residue in the process is linked to the emulsion on the film. The emulsion on the film is linked to the type of film and the process by which the film is developed.\nStandard films have an emulsion thereon which is compatible with the developing process for standard films which requires bleaching steps to convert the negative image on the film to a positive photographic image. Other films, such as direct positive films, do not require bleaching steps because a positive image is recorded directly onto the film when the film is exposed to that image. Photographic emulsions are provided on direct positive films which are compatible with this type of developing process.\nThus, there continues to be a need for a composition which will solubilize the oily dye residues produced by the development of direct positive silver halide film without foaming. If the oily dye residues are sufficiently solubilized, they will remain in solution and not adhere to the film or build up on the components of the film processor.\nIn continuous, automatic developers where the photodeveloping solution is recirculated and not replenished until a predetermined amount of film is developed, it is important that the oily dye residues remain solubilized until the photodeveloping solution is spent. Otherwise the film processor has to be shut down more frequently to change the solution, and the economy of the continuous, nonreplenished process is adversely effected. The composition also must not adversely affect the direct positive silver halide film developing process."}
{"text": "The present invention relates to disposable containers for hospital sharps and infectious waste, and pertains particularly to a secure disposable container and alternate closure assemblies for stand alone and for insert containers.\nHospitals and medical clinics use great quantities of sharps, such as needles, syringes, surgical blades, and the like, that are disposed of rather than cleaned and reused. It is necessary that the sharps be disposed of in a manner that prevents them being reused without sterilization. In particular, it is necessary to keep them from falling into the hands of those, such as intravenous drug users and the like, who are likely to use them without proper sterilization.\nNumerous containers have been developed in recent years, which are reasonably secure and disposable for receiving and disposing of hospital sharps, wastes and the like. Many of these containers are disposable inserts that mount into a security housing which is lockable and is securely mounted on a wall or other support structure. Many of these insert containers however have an inadequate security closure built into it, and are out dated so that they do not provide adequate security against pilfering of used syringes and the like therefrom. While improved containers have been developed which cannot readily be reopened and articles cannot be easily removed therefrom, such containers must be kept in a secure place or securely mounted to non-removable structure to prevent unauthorized removal.\nMany hospitals and clinics that have the older secure housings with insert type disposable containers do not wish to invest in new housings. It is therefore necessary to develop improved insert type disposable containers for these housings with improved closures for the combination. It is also desirable that new stand alone units be available.\nAn example of a prior art disposable insert type container and mounting housing is disclosed in U.S. Pat. No. 4,715,498 of Hanifl dated Dec. 29, 1987. This patent discloses a sharps container having a closure door that provides limited access to the interior of the container for mounting in a protective housing. However, an improved closure is desirable for such containers.\nExamples of stand alone type containers are disclosed in U.S. Pat. No. 4,736,860, granted Apr. 12, 1988 to Bemis, and entitled SHARPS DISPOSAL CONTAINER, and in U.S. Pat. No. 4,828,107, granted May 9, 1989 to Spencer and entitled DISPOSABLE CONTAINER FOR SYRINGES. It is desirable that simpler more effective disposable containers be available.\nIn the manufacture of multiple related or similar products, it is desirable that as many parts be interchangeable as possible in order to reduce manufacturing costs. It is also desirable that different variations of the product be available to meet different conditions and needs. It is, therefore, desirable that an improved securable disposable container combination be available to provide both insert and stand alone containers."}
{"text": "A single crystal used as a substrate of semiconductor devices, for example, a silicon single crystal, is mainly produced by Czochralski method (hereafter abbreviated to CZ method).\nWhen producing a single crystal by CZ method, for example, a single crystal production apparatus 10 as shown FIG. 2 is used to produce the single crystal. This single crystal production apparatus 10 has a member for containing and melting a raw material polycrystal such as silicon, heat insulation members to insulate heat, and etc. They are installed in a main chamber 11. A pulling chamber 12 extending upwardly is continuously provided from a ceiling portion of the main chamber 11, and a mechanism for pulling a single crystal 13 by a wire 14 (not shown) is provided above it.\nIn the main chamber 11, a quartz crucible 16 for containing a melted raw material melt 15 and a graphite crucible 17 supporting the quartz crucible 16 are provided, and these crucibles 16 and 17 are supported by a shaft 18 so that they can be rotated and moved upwardly or downwardly by a driving mechanism (not shown). To compensate for decline of the melt level of the raw material melt 15 caused by pulling of the single crystal 13, the driving mechanism for the crucibles 16 and 17 is designed to rise the crucibles 16 and 17 as much as the melt level declines.\nAnd, a graphite heater 19 for melting the raw material is provided so as to surround the crucibles 16 and 17. A heat insulating member 20 is provided outside the graphite heater 19 so as to surround it in order to prevent that the heat from the graphite heater 19 is directly radiated on the main chamber 11.\nMoreover, a graphite cylinder 23 is provided above the crucibles, and a heat insulating material 24 is provided on the outside of the lower end of the graphite cylinder 23 so as to oppose to the raw material melt 15 so that the heat radiation from the melt surface is intercepted and the temperature of the raw material melt surface is kept.\nA raw material lump is put in the quartz crucible 16 provided in the single crystal production apparatus as described above, the crucible 16 is heated by the graphite heater 19 as described above to melt the raw material lump in the quartz crucible 16. A seed crystal 22 fixed by a seed holder 21 connected with the lower end of the wire 14 is immersed into the raw material melt 15 melted from the raw material lump as described above. Thereafter, the single crystal 13 having a desired diameter and quality is grown under the seed crystal 22 by rotating and pulling the seed crystal 22. In this case, after bringing the seed crystal 22 into contact with the raw material melt 15, so-called necking, once forming the neck portion by narrowing the diameter to about 3 mm, is performed, and then, a dislocation-free crystal is pulled by spreading to a desired diameter.\nThe silicon single crystal produced by the CZ method as described above is mainly used for producing a semiconductor devices. In these years, with higher integration of semiconductor devices, an element becomes finer. The problem of Grown-in crystal defects introduced during the crystal growth has become more significant as an element has become finer.\nHereafter, the Grown-in crystal defects will be described (see, FIG. 4).\nIn the silicon single crystal, if the crystal growth rate is relatively high, Grown-in defects such as FPD (Flow Pattern Defect) considered to be originated from voids aggregated of void-type point defects, are present at high density over the region in a crystal radial direction. The region where the defects are present is called V (Vacancy) region. Moreover, if the growth rate is lower, OSF (Oxidation Induced Stacking Fault) region is generated as in a ring shape from the periphery of the crystal with lowering of the growth rate. And defects such as LSEPD (Large Secco Etch Pit Defect) and LFPD (Large Flow Pattern Defect) considered to be originated from dislocation loops aggregated of interstitial silicons in the outside of the ring are present at low density, and the region where the defects are present is called I (Interstitial) region. With further lowering of the growth rate, OSF ring is contracted to disappear in the center of a wafer, and the whole plane becomes I region.\nIn these years, in the outside of the OSF ring between V region and I region, a region has been discovered where neither such as FPD originated from voids nor such as LSEPD and LFPD originated from interstitial silicons are present. The region is called N (Neutral) region. Moreover, the N region is further categorized as follows. There are Nv region (a region where more voids are present) next to the outside of the OSF ring and Ni region (a region where more interstitial silicons are present) next to I region. In the Nv region, amount of precipitated oxygen is rich when thermal oxidation treatment is performed, and in the Ni region, amount of precipitated oxygen is little.\nThe Grown-in defects are considered to be determined the introduced amount by the parameter of F/G which is ratio of the pulling rate (F) and an average value of a temperature gradient in the crystal along a pulling axis from the melting point of silicon to 1400° C. (G) (See, for example, V. V. Voronkov, Journal of Crystal Growth, 59 (1982) 625-643). Namely, if the pulling rate and the temperature gradient are controlled so that F/G is constant, a single crystal can be pulled in a desired defect region or in a desired defect-free region (See, for example, Japanese Patent Laid-open (Kokai) Publication No. 2000-178099).\nTherefore, it is conventionally necessary to pull a single crystal by controlling the pulling rate etc. for N region so as to obtain a single crystal of defect-free region. It is difficult to control the rate because the single crystal of N region can be grown in a range of a relatively limited pulling rate. Thus, productivity and yield of the single crystal become low. Accordingly, a method of expanding a pulling rate range where a defect-free region can be obtained has been required so as to produce it more easily."}
{"text": "In ordinary metal microstructure inspection, common display methods of the microstructures after the microstructures are polished include an optical method, an etching method and an interference layer method, while the follow-up inspection of the microstructure and the analysis of metallography are mainly by means of the image observation and analysis of the microscope. The solution proportioning and reaction time of the etching method have a direct influence on the authenticity and clarity of displaying the metal microstructure, while image observation is limited by the amplification factor of the microscope. Therefore, the result of metallography has certain subjectivity and volatility, and cannot be accurately and quantitatively determined.\nIn the metal industry, the isostatic pressing technology is widely used to further improve the uniformity of the alloy microstructure. With the difference between the isostatic pressure and the mechanical property of the alloy microstructure, collapse and deformation will be different, meanwhile, the deformation behavior on the surface will be severe than the deformation behavior inside the alloy sample. Therefore, the isostatic pressing technology can be used to distinguish the microstructures with different hardness and mechanical properties on the surface of the alloy sample.\nWhen the isostatic pressure is used to carry out the microstructure metallographic determination with the fluid micro-explored strain, the microstructure can only be subjected to elastic plastic deformation in the vertical loading direction, and the microstructure will sink into a ‘horizontal potential well’ in the horizontal surface under the uniform extrusion effect of the pressure transmission medium and will not generate horizontal dislocation and slip. Therefore, no interference will be generated on the horizontal position of the microstructure.\nIn traditional metal microstructure inspection method, metallographic analysis is a macroscopic method, and rely more on subjective judgment and image magnification factor, while in the present invention, a technical method is adopted to characterize multiple parameters, when the full-field metallography of alloy microstructure is given, the statistical results of quantitative metallographic analysis can also be given."}
{"text": "The cellular telephone industry has had an enormous development in the world in the past decades. From the initial analog systems, such as those defined by the standards AMPS (Advanced Mobile Phone System) and NMT (Nordic Mobile Telephone), the development has during recent years been almost exclusively focused on standards for digital solutions for cellular radio network systems, such as D-AMPS (e.g., as specified in EIA/TIA-IS-54-B and IS-136) and GSM (Global System for Mobile Communications). Currently, the cellular technology is entering the so called 3rd generation (3G) by means of communication systems such as WCDMA, providing several advantages over the 2nd generation digital systems referred to above.\nMany of the advances made in mobile phone technology are related to functional features, such as better displays, more efficient and longer lasting batteries, and means for generating polyphonic ring signals. One functional feature which has become more and more common is built-in cameras. Cameras with video camera functionality are available today in several mobile phones. With the entrance of high bit-rate services, such as EDGE (Enhanced Data-rates for GSM) and 3G, the availability and usability for video-related services will increase. In particular, mobile video telephony, with simultaneous communication of sound and moving images, has recently become commercially available.\nFor stationary use, video conference systems generally include a camera mounted on or beside a communication terminal, such as a personal computer (PC), or integrated in an internet protocol (IP) enabled phone. Use of such a system may be fairly straightforward, as the user is positioned in front of the terminal with the camera aiming towards the user. However, mobile video conferencing is a bit more cumbersome. The terminal may be positioned in a support unit on a desktop, from which a camera in the unit is aimed towards the object of interest to be captured, typically the user. A more common way of using a mobile phone for video conferencing with face to face transmission is when it is held, such that the built-in camera is manually aimed towards the user. When communicating through a mobile handheld terminal, the user therefore may hold the terminal steady in front of the face so that the receiving party can see the face of the user, i.e. the sending party.\nA problem related to video conferencing with a radio terminal is caused by the fact that the built-in camera typically is placed adjacent to and parallel with the display, i.e. the optical axis of the camera is perpendicular to the display surface. The terminal therefore has to be aimed more or less 90° to the face, in order to get a proper image of the user. However, many users find this way of holding the terminal uncomfortable. Furthermore, for most mobile phone designs it may be difficult to use the terminal when placed on a desktop without additional supporting means, since it may require that the user's face be held over the terminal. A related problem is that the terminal may also include a small lamp aimed parallel with the camera to provide light to the object to be captured. When the camera, and the lamp, is aimed towards the face of the user at a 90° angle, there is also a risk that reflections of the user's face in the display surface will disturb the images presented on the display."}
{"text": "1. Field\nThe present invention relates to a metal capacitor and a manufacturing method thereof, and more particularly, to a metal capacitor in which an electric conductivity is significantly improved by applying a metal material for an electrolyte and a manufacturing method thereof.\n2. Background\nAn aluminum electrolytic capacitor is used to smooth a power output from a power circuit to be a predetermined value, or is used as a low frequency bypass. Hereinafter, a method of manufacturing the aluminum electrolytic capacitor will be briefly described.\nAn etching process of etching the surface of an aluminum foil is performed to enlarge a surface area of the aluminum foil and thereby increase an electric capacity. When the etching process is completed, a forming process of forming a dielectric substance on the aluminum foil is performed. When cathode and anode aluminum foils are manufactured through the etching process and the forming process, a slitting process of cutting the manufactured aluminum foil and a separator by as long as a desired width based on the length of a product is performed. When the slitting process is completed, a stitching process of stitching an aluminum lead patch, which is a lead terminal, to the aluminum foil is performed.\nWhen the slitting of the aluminum foil and the separator is completed, a winding process of disposing the separator between the anode aluminum foil and the cathode aluminum foil, and then winding the separator and the aluminum foils in a cylindrical shape and attaching a tape thereto, so as to not be unwounded. When the winding process is completed, an impregnation process of inserting the wound device into an aluminum case and injecting an electrolyte is performed. When the injecting of the electrolyte is completed, a curing process of sealing the aluminum case using a sealing material is performed. When the curling process is completed, an aging process of restoring a damage to the dielectric substance is performed. Through this, the assembly of the aluminum electrolytic capacitor is completed.\nDue to the current development in digitalization and thinness of electronic devices, when applying the conventional aluminum electrolytic capacitor, there are some problems as follow.\nSince the aluminum electrolytic capacitor uses the electrolyte, an electric conductive is comparatively low and thus a lifespan of the aluminum electrolytic capacitor is reduced in a high frequency area. Also, there are some constraints on improvement of reliability, a high frequency response, a low equivalent series resistance (ESR), and impedance. Also, due to a comparatively high ripple pyrexia, there are some constraints on stability and environments, such as fuming and firing."}
{"text": "1. Field of the Invention\nThis invention relates to novel fluorescent compositions and low-velocity electron excited fluorescent display devices utilizing the same, and more particularly is concerned with novel fluorescent compositions which can display emission having high luminance in a high state of color purity under low-velocity electron excitation, and low-velocity electron excited fluorescent display devices containing as a fluorescent screen these fluorescent compositions.\n2. Description of the Prior Art\nRecently, a low-velocity electron excited fluorescent display device (hereinafter referred to as \"fluorescent display device\") have been widely employed as a display device for desk top electronic calculators and various kinds of measuring instruments. As is well known, the fluorescent device of this kind in general has a fundamental structure such that both an anodic plate having a fluorescent screen on one side thereof and a cathode standing face to face with the above-described fluorescent screen are enclosed in an evacuated tube wherein the fluorescent screen placed on the anodic plate is excited by low-velocity electrons emitting from the cathode to result in emission of light of certain wavelengths. Both FIGS. 1 and 2 give outlines of typical structures of fluorescent display devices, and they show a diode type display tube and a triode type display tube, respectively. As shown in both FIG. 1 and FIG. 2, one side of an anodic plate 11 made of, e.g., an aluminium plate, has a fluorescent screen 12 thereon. The other side of the anodic plate 11 is supported by a ceramic base plate 13. The diode type display tube is equipped with a cathode standing face to face with the above-described fluorescent screen 12 placed on the one side of the anodic plate 11, and emission occurs by excitation of the fluorescent screen 12 which arises from low-velocity electrons emitted from the cathode 14. In particular, the triode type display tube shown in FIG. 2 additionally has a grid electrode 15 between the cathode 14 and the fluorescent screen 12 so as to control or diverge low-velocity electrons emitted from the cathode 14. Moreover, when the surface of the fluorescent screen 12 has wide area, two or more cathodes may be additionally placed in both fluorescent display tubes shown in FIG. 1 and FIG. 2 wherein only one cathode is placed, and there is no particular limit to the number of cathodes that can be placed therein. The aforesaid anodic plate 11 having a fluorescent screen 12 on one side thereof, the ceramic base plate 13 and the cathode 14 (which are shown in FIG. 1), or the aforesaid anodic plate 11 having a fluorescent screen 12 on one side thereof, the ceramic base plate 13, the cathode 14 and the grid electrode 15 (which are shown in FIG. 2) are enclosed in a transparent container 16, made of, for example, glass, the pressure inside which is held at a high vacuum of 10.sup.-5 to 10.sup.-9 Torr.\nZinc activated zinc oxide phosphors (ZnO:Zn) have been commonly known as phosphors employed for the above-described fluorescent display devices which can emit light of high luminance under low-velocity electron excitation occurring under certain conditions, particularly under acceleration potential below 100V. Phosphors of this kind can be prepared by firing zinc oxide (ZnO) alone in a reducing atmosphere, or by firing ZnO contaminated with a slight amount of a certain zinc compound other than ZnO such as zinc sulfide (ZnS) or the like in air, and they can give forth greenish white emission of high luminance when excited by low-velocity electrons. Fluorescent display devices having the fluorescent screen made of the aforesaid (ZnO:Zn) have been commercially used as display devices for, e.g., desk top electronic calculators and various kinds of measuring instruments. However, aside from (ZnO:Zn) almost no phosphors are known which can emit light under low-velocity electron excitation, and therefore, fluorescent display devices equipped with fluorescent screens containing phosphors other than (ZnO:Zn) are rare at the present stage of this art. Emission color of (ZnO:Zn) is greenish white as described above, and therefore a fluorescent display device utilizing (ZnO:Zn) has inadequate color purity as a green emitting display. Accordingly, the present invention is aimed at providing green emitting compositions and fluorescent display devices using them which can emit green light in a high state of color purity. In addition, the present invention is also aimed at providing fluorescent compositions capable of emitting light of specific wavelengths, other than greenish ones, in high luminance and at providing fluorescent display devices utilizing them."}
{"text": "1. Field of the Invention\nThe present invention relates to protective knee and elbow protectors employed by athletes and the like, particularly adapted to prevent injuries to the knee and the proximal anatomical structure associated therewith.\n2. Description of the Prior Art\nVarious knee and elbow protector devices are known in the art. It has been found that by placing a protector element atop the kneecap, injuries sustained in physical activity are reduced. Further, it has been determined by using such devices, the incidence of aggravating a current injury is also reduced. As such, knee and elbow protector devices are beneficial in the sport and athletic arts.\nExamples of knee and elbow protector protectors and other protector devices for use during athletic activity are described in the following U.S. Patents:\nU.S. Pat. No. 2,641,761 to Schultz discloses a knee brace or stabilizer. PA1 U.S. Pat. No. 3,044,075 to Rawlings discloses a protective device to protect the wearer at impact points. PA1 U.S. Pat. No. 3,322,873 to Hitchcock discloses a resilient body protector for use over various portions of the body. PA1 U.S. Pat. No. 3,670,725 to Gaylord discloses a protective pad to be employed generally with jointed portions of the body. This pad would be placed atop the jointed element to prevent debilitating bedsores and ulceration. PA1 U.S. Pat. No. 4,250,570 to Barlow discloses a protective knee support. PA1 U.S. Pat. No. 4,287,885 to to Applegate discloses a knee brace with a resilient pad for surrounding the patella. There exists no central pad, merely a central opening to receive the extended portion of the patella therethrough. PA1 U.S. Pat. No. 4,354,280 Hayes discloses a contact sport joint protector. PA1 U.S. Pat. No. 4,484,361 to Leighton discloses a knee and elbow pad and a method of manufacturing the same. PA1 U.S. Pat. No. 4,573,216 to Wortbeg discloses an impact dissipater to protect bones from impacts or falls. This device is primarily designed to be used proximal the bony region of the hip portion of the body. PA1 U.S. Pat. No. 4,700,698 to Aleylein discloses a knee orthosis. This device includes a right and a left lateral extension which is integral to the elastic cylindrical element. The device does not include a central foam element nor does it include an upper or lower lateral extension. PA1 U.S. Pat. No. 4,723,322 to Shelby discloses a knee pad. PA1 U.S. Pat. No. 4,796,303 Atwater discloses a knee protector pad. PA1 U.S. Pat. No. 4,870,956 to Fatool discloses a knee brace. PA1 U.S. Pat. No. 5,077,837 to Meistrell discloses a knee or elbow protector. PA1 U.S. Pat. No. 5,105,473 to Valtakwi discloses an athletic outfit having pockets for protective pads. PA1 U.S. Pat. No. 5,306,229 to Brandt discloses an elastic articular bandage. PA1 U.S. Pat. No. 5,411,037 to Hess discloses a knee-joint bandage. PA1 U.S. Pat. No. 5,474,524 to Carg discloses a joint support. PA1 U.S. Pat. No. 5,680,657 to Valtakwi discloses protective pads for insertion into pockets to protect from impact while playing sports. PA1 U.S. Pat. No. 5,685,021 to Tsujino discloses protective head gear for wrestlers containing a gel cushioning material. PA1 U.S. Pat. No. 5,695,452 to Grim discloses a formed resilient orthopedic pad. PA1 U.S. Pat. No. 5,769,809 to Witzell discloses a below the joint amputation limb protector. PA1 U.S. Pat. No. 327,961 to Decker, U.S. Pat. No. 341,005 to Pratt, U.S. Pat. No. 341,229 to Tsuiino, U.S. Pat. Nos. 360,284, 360,285, 360,286 and 360,287 to Pafet, U.S. Pat. No. 363,147 to Chen, U.S. Pat. Nos. 394,332 and 394,333 to Tsujino.\nThe following US design patents include pads and protectors for knees and other joints:\nAll of the forementioned patents are directed to protecting the knee or other body part from a direct impact. The instant invention is directed to a knee and elbow protector having a specific structure that protects the knee from an oblique, non-direct, or proximal impact. None of these references, either singly or in combination, disclose the instant invention and the benefits derived therefrom."}
{"text": "1. Technical Field\nThe invention relates to a multiple-ratio transmission mechanism in a powertrain for an automotive vehicle and to a control strategy for achieving smooth engagement and release of friction torque establishing elements during a transmission upshift event.\n2. Background\nIn a geared automatic transmission in an automotive vehicle powertrain having an engine or other torque source, a ratio change may be made from a so-called low ratio to a so-called higher ratio when a friction torque establishing element, such as a clutch or brake, is engaged in synchronism with disengagement of a companion friction torque establishing element. This is referred to as a ratio upshift. The friction torque establishing elements involved in the upshift may be referred to as an oncoming clutch or brake and an off-going clutch or brake. The upshift event is characterized by a preparatory phase, a torque phase and an inertia phase as the vehicle accelerates from a standing start.\nIn a conventional automatic transmission in a vehicle powertrain, the oncoming clutch torque capacity is controlled to increase from a low value during the torque phase. Simultaneous engagement of one clutch or brake and release of another results in a momentary activation of two torque flow paths through the gearing, causing a gear tie-up in which transmission output shaft torque decreases momentarily. This condition may be referred to as a “torque hole”. It occurs before the off-going clutch totally disengages.\nFriction elements, such as disc clutches, band brakes and disc brakes, typically are actuated hydraulically under the control of a transmission control module, which disengages an off-going friction clutch or brake while simultaneously engaging an oncoming friction clutch or brake during an upshift in order to lower speed ratio. For purposes of the present description of the invention, the clutch and the brake will be referred to as friction elements.\nDuring the preparatory phase, an automatic transmission control reduces off-going friction element torque capacity to prepare it for release as an actuator for the oncoming friction element is adjusted to prepare for its engagement. During the torque phase, the controller increases oncoming friction element torque capacity, which causes torque transmitted through the off-going friction element to drop quickly due to the transient gear tie up.\nAs torque is transmitted through the off-going friction element deceases, the automatic transmission output shaft torque drops, which causes the so-called torque hole. This is perceived by a vehicle occupant as an unpleasant shift shock. The inertia phase begins when the off-going clutch is released with no significant torque capacity."}
{"text": "1. Field of the Invention\nThe present invention relates to automotive vehicle glare reducing shields for use on vehicle sun visors.\n2. Description of the Prior Art\nIn the past, various sun visor extensions have been devised to serve as transparent extensions of an automotive vehicle sun visor. Such vehicle sun visors are conventionally foldably located near the upper boundary of a vehicle windshield. Sun visors of conventional design are opaque and are intended to serve as an optical shield at the upper portion of a windshield to prevent solar or other reflective glare from impairing a vehicle driver's vision. While the opaque sun visors perform this function under some glare conditions, many situations arise in which glare or high light level is directed through the windshield beneath the visor to nevertheless impair a driver's visual observation of road and traffic conditions. While conventional transparent shields have been provided and fixed to a vehicle sun visor, such devices of conventional design have had certain limitations.\nMany conventional transparent sun visor extension glare shields are not adapted to be folded flat into juxtaposition against an automotive vehicle sun visor when not in use. Rather, they have been affixed to the vehicle sun visor as rigid, readily noticeable, unsightly extensions thereof. Also, while useful under high glare conditions, such fixed sun visor extensions impair vision when glare is not a problem.\nOther glare shield extensions which have been designed for rotation about a bracket secured to the sun visor for storage purposes have not operated satisfactorily. Many such devices are constructed with plastic hinges or knuckles which are secured by friction to an axle or rod about which they are rotated. While the force of friction of the hinge knuckles on the horizontal mounting rod will hold the glare reducing sheet folded up in a stored position the first few times it is used, the plastic hinge knuckles rapidly tend to partially lose their elasticity, so that insufficient frictional force is present between the plastic hinge knuckles and the mounting rod to hold the glare reducing sheet upright in a position stored behind the sun visor. As a result, the glare reducing shield tends to flop downward and oftentimes assumes a position somewhere between its fully extended downwardly depending position, and its fully folded position behind the sun visor. This is both annoying and dangerous to the vehicle operator, and is especially irritating when the vehicle is traveling over rough and bumpy roads. Under such conditions the visor extension tends to tilt from behind the sun visor with considerable frequency. The vehicle operator's attention is thereupon distracted to a significant degree while the vehicle is traveling over surfaces which especially warrant the operator's undivided attention.\nVarious attempts have been made to adjust the diameter of the visor extension hinge pin and the plastic knuckle configuration to achieve a fit which is tight enough to allow the glare reducing shield to be positioned as desired and held in place by the force of friction between the plastic knuckles and the hinge pin. However, in each case the plastic knuckles deform inelastically over a relatively short period of time so that a tight friction fit between the plastic knuckles and the hinge pin cannot be maintained."}
{"text": "This invention relates to means for protecting marine installations such as offshore drilling and production installations in polar waters from large floating ice formations.\nMore specifically it relates to an energy absorbing system, which converts the kinetic energy of large ice structures into heat and power; the power being stored and used to energise, monitor and control, so that difficulties of control can be overcome.\nA number of ways exist for absorbing the kinetic energy of ice formations in polar waters. One way is to construct a gravel island on the seabed, with a long sloping beach, up which the ice formations ground.\nThe gravel island can be improved by containing the top section within a heavy caisson so as to minimise the quantity of filling material.\nAnother method consists of fabricating a number of large, hollow concrete horizontal slabs which can be floated into position and sunk onto the seabed stacked on top of each other. The finished face of the outer wall is shaped so that the advancing ice is forced upward until the horizontal motion becomes vertical and is absorbed by gravity.\nA fourth method is to create underwater obstacles, sometimes in the form of a vertical wire fence, anchored to the seabed, against which the advancing ice keels will collide.\nA fifth method is by creating an underwater circle of drag elements (as described in my U.S. Pat. No. 4,547,093 of Oct. 15, 1985) which are attached to seabed anchorages by long anchor lines, the drag elements having teeth which penetrate the underside of the ice and absorb the kinetic energy of the ice formation by their shearing action.\nAll these methods suffer from very real disadvantages. The artificial islands are very expensive and take a long time to construct even in shallow water. Removal is an even greater problem. The caisson retained top helps to minimise the cost, but there are still time and cost problems. Concrete slabs stacked on top of each other have the advantage of re-use on a number of different sites in different seat water depths, but they become very large and expensive in deep water especially if they have to resist the collision impact force from a large ice island. Sea obstructions can be destroyed too easily by a large ice keel, and repair is not possible during the winter season. Drag elements can be very effective with fast moving ice but they suffer from the disadvantages that the drag force cannot be controlled, the teeth-freeze into very slow moving ice, and there can be a tendency for following pack ice to overslide ice slabs held on the drag elements."}
{"text": "A video telephony service enables a user to communicate while viewing images of the other party through a mobile communication terminal to which a camera is attached, and can be classified into a circuit-switched method using a Mobile Switching Center (MSC) and a packet-switched method using an all-Internet Protocol (IP) network.\nFIG. 1 is a block diagram showing an example of the construction of a mobile communication network for providing video telephony service.\nWhen a video telephony call is attempted by a calling mobile communication terminal 10, a base station 12 receives the call, and performs video telephony with a called mobile communication terminal 20 through an MSC 26, a Base Station Controller (BSC) 24, and a base station 22 that are connected with the called mobile communication terminal 20 through a BSC 14 and an MSC 16. Video telephony service, as described above, is provided through the MSCs 16 and 26 in the case where the circuit-switched method is employed, and the service is provided through a General Packet Radio Service (GPRS) providing register, that is, a Serving GPRS Support Node (SGSN) 30, and a GPRS gate providing register, that is, a Gateway GPRS Support Register (GGSN), that are connected with the BSCs 14 and 24, in the case where the packet-switched method using an all-IP network is employed.\nIn order to provide video telephony service, an international standard protocol, for example, H.323 or H.324M, defined by International Telecommunication Union (ITU), must be implemented in a mobile communication system. H.323 is a system protocol that enables the provision of video telephony service in an IP network, that is, a packet data network, H.324 is a system protocol developed on the basis of a public network, and H.324 Mobile (M) is a system protocol that has been improved for mobile communication.\nIn a video telephony service using H.324M as a system protocol, the compression and encoding of moving images are performed using H.261 and H.263, and the encoding of voice is performed using G.723.1. With regard to this, H.261 is a moving image compression/encoding standard for video telephony and video conferencing, and H.263 and MPEG-4 are video compression/encoding standards that have been improved more than H.261. Furthermore, G.723.1 is a standard for converting voice signals to be less than 8 Kbps. A video standard used in 3GPP is 3G-324M, which was created by modifying H.324M to be suitable for 3GPP. 3G-324M is considerably different from H.324M in that Adaptive Multi-Rate (AMR) is basically used as a voice codec, and G.723.1 is an option.\nFurthermore, H.324M uses H.223 to multiplex moving images, voice and data, and uses H.245 to allocate a point-to-point or point-to-multipoint channel by selecting a voice codec and performing a function of logical channel signaling.\nFIG. 2 is a flowchart illustrating a general video telephony service method, and shows an example of using a 3G-324M protocol.\nAs shown in the drawing, when a calling mobile communication terminal attempts a call to use video telephony service, a setup process is performed using H.223 at step S10 and, thereafter, the call is set up at step S20.\nAfter the call setup has been completed, a line connection and allocation process based on negotiation between calling and called mobile communication terminals using H.245 is performed at step S30. After a communication path has been established accordingly, moving images, voice, and data (picture, photograph and the like) are exchanged between the calling and called mobile communication terminals at steps S40 to S60.\nFIG. 3 is a flowchart illustrating the line connection and allocation process of FIG. 2 in detail.\nThe line connection and allocation process is a process of exchanging information about the characteristics of mobile communication terminals and, thereby, making settings so that video telephony is made possible, in the case where voice codecs or video codecs used for the calling and called mobile communication terminals are different from each other.\nAs shown in the drawing, for line connection and allocation between the calling and called mobile communication terminals, a master/slave decision and response process is performed at step S310, a terminal characteristic information exchange process at S320, a multiplexing information exchange process at step S330, a logical channel generation process for voice transmission at step S340, and a logical channel generation process for video transmission at step S350.\nAs described above, the current negotiation process for video telephony is complicated, and a lot of time is required, so that a problem occurs in that the communication path is not established immediately after call setup, and communication is not performed, and the delay time occurs in proportion to the time for which the negotiation process is performed."}
{"text": "1. Field of the Invention\nThe present invention relates to a data error detecting apparatus and a data error detecting method. The present invention particularly relates to a data error detecting apparatus and a data error detecting method in a single-chip microcomputer.\n2. Description of the Related Art\nSemiconductor devices called as microcomputers are installed in many commercially available apparatuses. The microcomputer has a microprocessor (CPU) for performing arithmetic and logic processing and a memory, which are generally formed as an integrated circuit (IC) on a single LSI chip.\nFor apparatuses with microcomputers installed, there are apparatuses that reliability is especially important. For example, electric components provided for vehicles require high reliability at all times. For this reason, the microcomputer provided to such an apparatus is designed based on a design concept called fail-safe, in order to measure unexpected trouble occurrence. Such a microcomputer is designed under the presumption that problems such as a failure, an operational mistake, and a design defect occur. In occurrence of the problem, the microcomputer has a function for keeping damage minimum. For example, when a CPU in a system reads out extremely important data, the fail-safe function instantaneously detects data abnormality resulting from device failure by always monitoring the data to prevent occurrence of false operations.\nOne fail-safe function is realized as a data error detecting apparatus in a microcomputer. A parity bit is added to important data, and a parity check is carried out when a CPU reads the data. If any abnormality is found, parity error interrupt is issued to inform data abnormality as soon as possible.\nFIG. 1 is a block diagram showing a configuration of a conventional microcomputer 101 having a data error detecting apparatus. With reference to FIG. 1, the conventional microcomputer 101 includes a peripheral function 102, a data error detecting apparatus 103, and a CPU 104. Additionally, as shown in FIG. 1, the data error detecting apparatus 103 includes a parity bit adding circuit 105, a register 106, and a parity check circuit 107, which are connected through a bus 109.\nIn the conventional microcomputer 101, the parity bit adding circuit 105 adds a parity bit to a data in storage of the data, and writes the data with the added parity bit into the register 106. The parity check circuit 107 carries out parity check when the CPU 104 reads out the data from the register 106. The parity check circuit 107 determines whether there is any breaking of a data line and failure in a circuit for controlling the data, based on the result of the parity check. When it is determined that there is breaking of the data line or a failure in a circuit for controlling the data, the parity check circuit 107 issues a parity error interrupt, thereby detecting data abnormality.\nAs stated above, the conventional data error detecting apparatus 103 detects a failure in the microcomputer 101 by carrying out parity check to detect a data error corresponding to the data line (data line and the circuit for controlling the data). In general, the data error detecting apparatus 103 does not monitor any failure in the control circuit for writing the data into the register, breaking of a clock line, and so on. Therefore, the data error detecting apparatus 103 cannot know whether or not the data is originally right. In other words, the data error detecting apparatus 103 regards a currently retained data as a correct data and continues a reading operation even when data update is no longer possible.\nFIG. 2 shows an operation of the data error detecting apparatus 103 in the conventional microcomputer 101 when there is occurrence of a failure in a control circuit writing data into a register or breaking of a clock line. With reference to FIG. 2, the following data are outputted from the peripheral function 102:\nFirst time: 0111;\nSecond time: 0101;\nThird time: 1110; and\nFourth time: 1010.\nAlso, it is shown that data update is not carried out due to the breaking of a line when the data should be updated from second data to third data. As shown in FIG. 2, the parity check circuit 107 reads non-update data when data update is not possible due to the breaking of the line. At this time, an error is not detected in the parity check since a consistency between a register value and a parity bit is kept. For this reason, the CPU reads an old data as a correct data, leading a delay to detection of trouble occurrence. Therefore, the conventional data error detecting apparatus 103 continues data readout without determination of a failure even if a read value is the same value for ten consecutive times."}
{"text": "Safety instrument systems typically incorporate emergency shutdown valves which are normally in a fully opened or a fully closed state and are controlled by a logic solver, a Programmable Logic Controller (PLC), or an emergency shutdown controller of some type to change states in the event of an emergency situation. To ensure that these valves can function properly, process control system operators typically periodically test the emergency shutdown valves by running these valves through a stroke test, which partially or completely opens or closes the valve. Because these tests are typically performed while the process is operating on-line or is operational, it is important to perform any test reliably and then return the valve to its normal state as quickly as possible. In this context, the term “normal state” refers to the position or state of the emergency shutdown valve when there is no emergency and the emergency shutdown valve is not being tested, i.e., when the process is operating normally.\nIn many cases, the emergency shutdown tests are performed at predetermined intervals by remotely located controllers. For example, emergency shutdown tests may be performed only a few times each year due to cumbersome test procedures and issues related to manpower. Also, during emergency shutdown tests, the emergency shutdown valve, or other emergency shutdown device being tested, is not available for use if an actual emergency event were to arise. However, limited, periodic testing is not an efficient way of verifying the operability of an emergency shutdown test system. As a result, digital valve controllers have been, in some cases, programmed to assist in the operation of the valve test to make the testing more automatic, user friendly and reliable.\nAdditionally, it is typically important that any emergency shutdown system be able to activate an emergency shutdown device (an emergency shutdown valve, for example) to its safe condition even when commanded by the emergency shutdown controller to do so in the unlikely but possible situation where an emergency event occurs during an emergency shutdown device test. In this context, the term “safe condition” refers to the position of the emergency shutdown device that makes the process plant or portion of the process plant “safe.” Typically, this safe position is associated with a position of the shutdown device that shuts down or halts some portion of the process plant.\nWhile there are many systems that test the ultimate emergency shutdown device, such as an emergency shutdown valve, itself, in many cases there is supporting equipment associated with the emergency shutdown device that should also be tested to assure the complete operability of the emergency shutdown capabilities at any particular plant location. For example, in some pneumatic valve configurations, a solenoid valve is connected between a pneumatic valve actuator of an emergency shutdown valve and an emergency shutdown controller to redundantly control the operation of the valve actuator in response to signals from the emergency shutdown controller. While the emergency shutdown valve may be functional, it is possible for the solenoid device to become defective and therefore not operate properly as a redundant method of actuating the emergency shutdown valve. In some cases, an improperly operating solenoid device may even prevent the emergency shutdown valve from actuating properly when the emergency shutdown controller sends a shut-down signal to the valve controller for the emergency shutdown valve.\nWhile it is possible to develop and provide specialized equipment at each emergency shutdown location within a plant to perform testing of each different emergency shutdown device and its supporting equipment, it is more desirable to provide a universal or generic set of equipment that may be used in many different situations to test different types of emergency shutdown devices and the supporting equipment associated therewith or to perform other functions in the plant. For example, it is desirable if such versatile equipment is able to control and test different types of emergency shutdown valves and solenoid valve configurations while simultaneously or alternatively operating as part of a closed loop distributed process control system."}
{"text": "The present invention relates to an apparatus for printing upon the wall of a container, whereby during printing the wall of the container is moved around a rotatable support which keeps the container wall pressed directly against a rotating printing surface.\nFor printing upon or painting containers it has hitherto been the practice to use an internal support whose circumference abuts precisely upon the inner wall of the container whereby it is possible with the aid of this rotatable support to carry out continuous printing either only upon circular beakers or only upon substantially rectangular containers of a specified size and shape."}
{"text": "As fixed type constant velocity universal joints, there have been publicly known joints of a six-ball Rzeppa type (BJ), a six-ball undercut-free type (UJ), an eight-ball Rzeppa type (EBJ), an eight-ball undercut-free type (EUJ), and a track groove crossing type in which paired track grooves cross each other (see, for example, Patent Literature 1). Those joints are used as appropriate in accordance with purposes, required characteristics, and the like.\nReferring to FIG. 14a and FIG. 14b, description is given of an example of the fixed type constant velocity universal joint of the track groove crossing type (hereinafter referred to as “crossing fixed type CVJ”). FIG. 14a is a vertical sectional view (sectional view taken along the line A-O-B in FIG. 14b) of a state in which the crossing fixed type CVJ for a propeller shaft forms an operating angle of 0°, and FIG. 14b is a front view of the crossing fixed type CVJ. The constant velocity universal joint 100 includes an outer joint member 102, an inner joint member 103, balls 104, and a cage 105. Eight arc-shaped track grooves 107 are formed in a spherical inner peripheral surface 106 of the outer joint member 102. The track grooves 107 are formed so that planes including ball raceway center lines x of the track grooves 107 are inclined with respect to a joint axial line n-n and the track grooves 107 are adjacent to each other in a circumferential direction with their inclination directions opposite to each other. Eight arc-shaped track grooves 109 are formed in a spherical outer peripheral surface 108 of the inner joint member 103. The track grooves 109 are formed so as to be mirror-image symmetrical with the paired track grooves 107 of the outer joint member 102 with respect to a plane P including a joint center O at the operating angle of 0°. That is, the inner joint member 103 is assembled to an inner periphery of the outer joint member 102 so that the paired track grooves 107 and 109 cross each other.\nAs illustrated in FIG. 14a, curvature centers of the track grooves 107 of the outer joint member 102 and the track grooves 109 of the inner joint member 103 are each positioned at the joint center O. The balls 104 are interposed in crossing portions between the paired track grooves 107 and 109, and are received and held in pocket portions 105a of the cage 105 arranged between the spherical inner peripheral surface 106 of the outer joint member 102 and the spherical outer peripheral surface 108 of the inner joint member 103. Curvature centers of a spherical outer peripheral surface 111 and a spherical inner peripheral surface 112 of the cage 105 are each positioned at the joint center O. In the constant velocity universal joint 100, the balls 104 are interposed in the crossing portions between the paired track grooves 107 and 109. Therefore, when the joint forms an operating angle, the balls 104 are always guided in a plane bisecting an angle formed between axial lines of the outer joint member 102 and the inner joint member 103. With this, rotational torque is transmitted at a constant velocity between the two axes.\nFurther, the track grooves 107 of the outer joint member 102 and the track grooves 109 of the inner joint member 103 are adjacent to each other in the circumferential direction with their inclination directions opposite to each other. Therefore, when both the joint members 102 and 103 rotate relative to each other at the operating angle of 0°, forces in the opposite directions are applied from the balls 104 to the pocket portions 105a of the cage 105 that are adjacent to each other in the circumferential direction. With this, the cage 105 is stabilized at the position of the joint center O. Thus, a contact force between the spherical outer peripheral surface 111 of the cage 105 and the spherical inner peripheral surface 106 of the outer joint member 102, and a contact force between the spherical inner peripheral surface 112 of the cage 105 and the spherical outer peripheral surface 108 of the inner joint member 103 are suppressed. As a result, it is possible to attain the constant velocity universal joint 100 that is capable of suppressing torque loss and heat generation, and is excellent in durability. Thus, with use of the constant velocity universal joint 100, it is possible to attain a propeller shaft that is suppressed in torque loss and heat generation, and is enhanced in efficiency.\nOperability of fixed type constant velocity universal joints including the constant velocity universal joint 100 described above, in which rotational torque is transmitted between both the joint members through intermediation of the balls, is secured by setting wedge angles of ball tracks formed by the paired track grooves (by holding the balls with the paired track grooves). In the constant velocity universal joint 100 of the track groove crossing type, as described above, when wedge angles, which are formed so that a force of pressing the cage 105 to one axial side (for example, right side in FIG. 14a) is applied from the balls 104 at the operating angle of 0°, are defined as positive wedge angles, and wedge angles, which are formed so that a force of pressing the cage 105 to the other axial side is applied from the balls 104 at the operating angle of 0°, are defined as negative wedge angles, ball tracks that form the positive wedge angles and ball tracks that form the negative wedge angles are alternately formed in the circumferential direction. With this, desired joint performance is secured. In the constant velocity universal joint 100 described above, the curvature centers of the track grooves 107 and 109 are each positioned at the joint center O, and hence sizes and opening directions of the wedge angles are determined in accordance with angles (inclination angles) that are formed by the track grooves 107 and 109 (by planes including the ball raceway center lines thereof) with respect to the joint axial line n-n, contact angles of the balls 104 with respect to the track grooves 107 and 109, and operating angles."}
{"text": "The present application discloses an access control apparatus and an access control method for controlling accesses to each recording area on a disk on a command with a time limit requiring completion of the command in a predefined period of time.\nRecently, magnetic disk apparatuses (hereinafter referred to as HDDs) are widely used as data storage incorporated in personal computers (hereinafter referred to as PCs) or connected to the PCs. The use of the HDDs has been expanding in application, not only for PCs, but for various equipment such as video recorders and car navigation systems requiring storage of a large amount of data.\nNot all sectors in a recording area of the magnetic disk (hereinafter referred to as an HD) mounted in the HDD are writable or readable. There are defective sectors in the recording area to which data can not written or from which data can not read. Therefore, the HDDs require a process which is called an automated alternate process.\nFIGS. 1A, 1B, and 1C are explanatory illustration of the automated alternate process.\nEach sector included in the recording area of the HD is assigned a logical block addressing (LBA) address for identifying the sector as shown in FIG. 1A. In this application, data are written onto each sector in numerical order as LBA 10, 11 and 12. For example, if writing data onto a sector at LBA address 13 fails and the data are not written when retrying the access a predefined number of times, the data attempted to be written in the sector at LBA address 13 are stored in a sector at LBA address 0 as shown in FIG. 1B. Then, information that the data attempted to be written in the sector at LBA address 13 are stored in the sector at LBA address 0 is recorded in an alternate candidate table shown in FIG. 1C provided in a different area of the HD.\nTherefore, the data are read from the sector at LBA address 0 instead of the sector at LBA address 13 when a read command to read the data stored in the sector at LBA address 13 is received.\nAs described above, the process for storing data on the alternate sector instead of the sector on which the data are initially attempted to be written is called an alternate process. An alternate process that is executed automatically in a specific sequence is called the automated alternate process.\nHowever, the automated alternate process takes a long time when an access to a sector is retried and failed. The automated alternate process takes a long time because a magnetic head is moved to the alternate area to store the data and the information on the alternate sector is recorded in the alternate candidate table.\nFor example, stream commands such as movies and music require real-time processing. These stream commands are issued with a command completion time limit (CCTL). The command completion time limit requires execution of the commands in a period of time specified by the time limit absolutely. If the accesses to multiple sectors are retried during execution of the stream command, a time-out error will occur and the movie or the music may be intermittent.\nJapanese Laid-open Patent Publication No. H11-327808 discloses that the automated alternate process is not executed when the time is running out even though the process is required.\nIn Japanese Laid-open Patent Publication No. H11-327808, the automated alternate process is not executed. Therefore, data to be stored at the sector are not stored. Thus, the data to be read from the sector does not exist. This leads to deterioration of a data quality such as an image quality of the movie or a sound quality of the music.\nJapanese Laid-open Patent Publication No. 2003-196924 discloses that the automated alternate process is not executed immediately after failing to store the data. The data are stored in an interim buffer and the data are read from the buffer and stored on the alternate sector of the HD when processing time is finished.\nHowever, Japanese Laid-open Patent Publication No. 2003-196924 discloses only the write process onto the HD, and mentions no comprehensive process including a read process. Besides, Japanese Laid-open Patent Publication No. 2003-196924 does not point out an issue or offer a solution when a write command is reissued and the sector at the LBA address is accessed again before the data stored in the interim buffer are stored on the HD or a read command is issued and the sector at the LBA address is accessed."}
{"text": "The present invention relates to a method of improving the control behavior of an anti-lock control system, wherein a vehicle reference speed that is used as a reference quantity to determine the wheel slip of the individual wheels is derived from the rotating behavior of the vehicle wheels, and criteria are obtained for identifying a cornering situation and the curve direction, wherein for identifying cornering the wheel slip values are filtered by way of a low-pass filter and the filtered values are compared.\nIt is already known in the art to extend the functions of an anti-lock control system by employing the system for improving the driving stability or deceleration stability in the curve. For this purpose, during cornering or during a partial deceleration, i.e., in a deceleration process in which the threshold of response to the anti-lock control system is not achieved, a stabilizing moment about the vertical axis of the vehicle is created by specifically delaying the brake pressure build-up on the wheels at the inner side of the curve compared to the brake pressure on the wheels on the outer side of the curve (`Bremsanlage und Schlupf-Regelsystem der neuen 7-er Reihe von BMW` (Brake System and Slip Control System of the new BMW No. 7 series); ATZ 97 (1995), pp. 8-15, and `Bremsanlage und Schlupf-Regelsysteme der neuen Baureihe 5 von BMW` (Brake System and Slip Control Systems of the new BMW No. 5 series); ATZ 98 (1996), pp. 188-194. The information on the actual steering angle is derived--if no steering angle sensor is used--from the transverse acceleration which, in turn, is calculated from the wheel sensor signals.\nDE 34 13 738 C2 discloses an anti-lock control system (ABS) provided with a system for identifying a cornering situation based on wheel slip measurement. To identify a cornering situation, the sum of slip values on the wheels at one vehicle side are summed up and compared to the slip values on the wheels at the other vehicle side, generating a cornering identification signal as soon as the difference of the summed-up slip values exceeds a predetermined limit value. Upon identification of a cornering situation, selective criteria, such as `select-low` or `select high` according to which the pressure variation in the individual brake pressure control channels of the said brake system is controlled, and limit values for rendering the said selective criteria effective, are varied. In this way, the control will be adjusted to the varying conditions during straight-forward driving and during cornering.\nDE 21 19 590 A1 which has an older filing date teaches generating a cornering identification signal with the aid of a transverse acceleration meter, such as a mercury switch.\nThe calculation of the filtered wheel slip and the evaluation of this quantity to determine criteria for cornering are described in the patent application DE 195 22 634 A1 which is not prior published. The filtered wheel slip values are compared and evaluated for cornering identification according to the above publications. Cornering is identified when simultaneously the values of the filtered wheel slip on both front wheels exceed a predetermined maximum slip value and the filtered wheel slip of a rear wheel is in excess of an equally predetermined maximum slip value and the filtered wheel slip of the second rear wheel is below a determined minimum slip value.\nWhen cornering is identified, switch-over to a special, i.e., cornering, control mode is effected according to patent application DE 195 22 632 A1 which is also not prior published.\nAn object of the present invention is to improve, in a method of the type referred to hereinabove, the ability of evaluating the measuring signals and, more particularly, the reliable identification of a cornering situation compared to straight travel."}
{"text": "1. Field of the Invention\nThe present invention relates generally to methods and systems using local number portability (LNP) databases and, more particularly, to a method and system for forecasting telephony traffic trends using LNP databases.\n2. Background Art\nLocal number portability (LNP) allows wireless and landline telephone customers to keep their telephone numbers when they change residences and/or service providers. A LNP database stores a list of telephone numbers that have been switched or ported from one service provider to another service provider. For each ported telephone number, the LNP database also stores an associated location routing number (LRN). The networks of the service providers use the LRN to route a telephone call from a point of origination (such as from a telephone associated with the ported telephone number or from another party's telephone) to a point of destination (such as to the other one of the telephone associated with the ported telephone number or to the other party's telephone) over the public switched telephone network (PSTN).\nPorting telephone numbers between service providers causes changes in the way that the networks of the service providers handle telephone calls associated with the ported telephone numbers. Generally, the network of a service provider handles more telephone calls if there is a net influx of ported telephone numbers to the service provider. As such, the service provider should ensure that its network has enough capacity and is properly connected to the networks of other service providers in order to handle the additional telephone traffic. Likewise, the network of a service provider handles less telephone calls if there is a net outflow of ported telephone numbers from the service provider. As such, this service provider should consider changing its network and reconfiguring how its network is connected to the networks of other service providers in view of the reduced amount of telephone traffic.\nThe networks of the other service providers not directly involved in the porting of telephone numbers between two service providers may also handle telephone calls associated with the ported telephone numbers differently after these telephone numbers have been ported between the two directly involved service providers. As such, these other service providers should consider changing their networks and reconfiguring how their networks connect with the networks of other service providers in view of the differently routed telephone traffic.\nPresently, network engineers monitor service provider networks using broad measures such as minutes of use on trunk groups, blocking on trunk groups, and point-to-point telephone traffic flows through the networks. These metrics are compiled over long periods of time (for instance, several months) and therefore are indicative of what has previously happened in the networks. As such, these are lagging indicators of telephone traffic trends.\nThus, there exists a need for using LNP database images or snapshots obtained at periodic times to track ported telephone numbers between service providers and then use this data along with topology and interconnection data regarding the networks of the service providers in order to predict future telephone traffic on the service provider networks."}
{"text": "1. Field of the Invention\nThe present invention relates generally to juvenile products, such as multiple infant activity centers for use by children and their caregivers, and more particularly to height adjustment mechanisms for use in connection with such products.\n2. Description of the Related Art\nJuvenile products are widely used by caregivers. Many of those products, such as multiple infant activity centers, have seating portions for a child, or other aspects, that occasionally need adjustment to accommodate children of different sizes. It is generally desired that such adjustment quick and easy for the caregiver, but result in a secure adjustment.\nWhile various devices have been used in the past to accomplish such adjustment, most of those prior art devices have been lacking in ease of adjustments, security, or both.\nIn sum, the prior art devices do not provide the important advantages of allowing easy adjustment of the juvenile products, such as seat position in multiple infant activity centers."}
{"text": "1. Field of the Invention\nThe present invention is related to methods of protection against malware located on web resources and, in particular, to malware scans of web resources and identification of malware components on web resources.\n2. Description of the Related Art\nDetection of viruses and malware has been a concern throughout the era of the personal computer. With the growth of communication networks such as the Internet and increasing interchange of data, including the rapid growth in the use of e-mail for communications, the infection of computers and networks through communications or file exchanges is an increasingly significant consideration. Infections take various forms, but are typically related to computer viruses, Trojan programs, or other forms of malicious code (i.e., malware).\nRecent incidents of e-mail mediated virus attacks have been dramatic both for the speed of propagation and for the extent of damage, with Internet service providers (ISPs) and companies suffering service problems and a loss of e-mail and networking capability. In many instances, attempts to adequately prevent file exchange or e-mail mediated infections significantly inconvenience computer users.\nPopularization of web services increases malware-related threats to web clients. With the development of web-technologies, such as AJAX, JAVA, PHP, FLASH, etc., web sites become more accessible, with more functionality with various media content, which enhances their appeal and popularity among users. With the introduction of methods of designing web-systems, known as Web 2.0, users have been enabled to not only receive information from the Internet, but also fill out web sites with their own data.\nMany Internet users have their own blogs, pages within social networks, where they can share messages or data with other users. In addition to text messaging, modern technologies support transfer of media files, such as photos, video files, interactive documents, animations and applications.\nActive participation of users in creation and modification of a web site leads to a rapid development and changes of the site content. A typical example of such a site is a news ticker, where new articles with links to news reports, photos or videos appear at intervals of a minute and sometimes even more often. Another example is a web forum, where a number of users can be in excess of several thousand, and where new messages appear every second.\nIncreasing popularity of web resources makes this environment more attractive to hackers and virus writers, who spread malicious programs over the Internet. A malicious script, such as iframe, exploits, etc. can be added to the files uploaded to the site. Thus, file exchanges, forums, blogs, web-interfaces of mail servers and any other resources can be infected.\nThe infection can be perpetrated through any interface, such as a usual web site form to be filled out by a user (in this case, the infection is performed mostly manually by an insider or an intruder who stole passwords). A perpetrator can exploit vulnerability in the content management system or file access the site via FTP. The list of resources is not limited to the http transfer protocol but also includes ftp-resources and other servers.\nTypically, ftp-servers are used as tools for remote administration of the sites, including editing and uploading scripts. However, a connection to the ftp-server provides the ability to test all scenarios of site, and not just those that are executed on the user's computer. Access the site files via FTP allows a security system to analyze the source scripts and original pages, while access via http allows to examining only the result of the script and/or the result of processing of the web site files by a web server.\nAn unauthorized access to a user account and uploading on behalf of this user, a malicious program could lead to a rapid spread of malware because of the credibility of that person among other users. When dealing with mundane sites, users' attention may fade over time due to the usual circumstances. Because of this, incidents of exploitations of social networks and phishing are increasing.\nIn addition to public Internet services, there are some local sites (corporate, local area networks, user group sites, etc.) that are not accessible from outside of their network. This imposes significant problems in malware scanning of a given resource by anti-virus services or laboratories.\nA typical scheme of interaction between a user computer system and a Web resource is illustrated in FIG. 1. Web-server is an application that performs the functions of the server (i.e., a computer system) on which the application is implemented. In addition to web-servers 110 or ftp servers 120, other server applications can be installed. For example, mail servers, proxy servers, IRC-server, etc. In order to interface with these servers, client (i.e., a user) computers 140 must have special applications installed, such as, a browser 130, a file manager 150, etc.\nClient 140 access the web server 100 using a URL address of a desired web page or other resource. Each server application and a client interact, using different protocols. The primary protocols for data transfer between the client 140 and server 100 are HTTP (HTTPS), (S) FTP, POP3, SMTP, IMAP4, etc.\nMost servers support authentication in the following manner: an authorization data is transferred from the client 140 to the server 100, where the authorization takes place. Then, the data is transmitted from the server 100 to the client 140 based on the rights granted to the client 140.\nPersonalization makes it possible to make a web resource unique to each user. Authorization is used to distinguish users, their rights and available data. User authorization is carried out by the user via a client application, for example, by filling the forms in the browser. Very often these applications allow preserving the identity of the user. Thus, each user can customize his own interface, to restrict access to his personal page or email, to identify the displayed pages by sections or topics.\nIn order to provide security while using Internet or on the local network, downloadable content needs to be checked. A downloaded page may contain viruses, trojans, adware, spam, exploits (i.e., HTML code, links to image or “.pdf” file, which has special modifications that result in browser errors and execution of some potentially harmful code) for applications (such as for example, PDF-reader, web browser, media player, flash-player, etc.)\nA system of protection of personal computers can include a file anti-virus (AV), network screen, a firewall, a special protection against network attacks, Web anti-virus as well as remote security means of AV vendor companies. Currently AV technologies are heavily developed and includes many different methods and systems that implement both heuristic and signature analysis.\nThe signature type web analysis includes:\nassembling black list of pages (URL-blacklist);\nassembling white lists of trusted (i.e., clean) applications/components; and\nstoring a collection of malware components.\nThe heuristic analysis typically includes any of:\nemulation of executable programs/components;\nemulation of executable scripts;\nvirtualization of execution environment; and\ncontrol and analysis of application activity.\nA firewall or a network screen is necessary for monitoring and filtering network packets at different levels of the OSI model in accordance with assigned rules for network connections: filtering based on static rules and filtering with tracking executed applications and controlling logic and algorithms of the relevant protocols.\nA system for protection against network attacks is typically launched at system startup and monitors incoming traffic activity for patterns typical of network attacks. If an attempt to attack a computer is detected, the system blocks any network activity for the attacker to a protected computer.\nA conventional web anti-virus (web AV) intercepts and blocks execution of a malicious script on a web site if it poses a threat. Strict control is also imposed over all HTTP-traffic. Web AV also analyzes web resources for phishing scams and filters banners and pop-ups.\nAnother line of defense can be a local security server that analyzes the situation within a local network, scans the local traffic and analyzes the network activity of computers. There are also web services that allow the verification of the Internet resource or files. The user downloads a file or enters a URL of the resource and the entire malware test takes place on the web-service of an AV company.\nIn the case of local AV remedies, a web resource is checked when it is downloaded, in other words, when the client application goes to the corresponding URL-address. It is important to note that in this case, the user is authenticated on the web resource and downloads content, as defined for his user account. The content can be in form of scripts, links, articles, messages, reports, letters from trusted users, etc.\nFIG. 2 illustrates a conventional system for checking a web resource 200 for presence of malware or links to infected resources. A connection with the resource 200 is established by an application 220 of a client 210. The client application 210 in this case can be a regular browser, a file manager or another application that interacts with a server 200 via data transfer protocol 230.\nThe client application 220 transmits to the server 200 user identifiers 250. Depending on authentication scheme 260, the identifiers 250 can represent logins, passwords, session keys, cookie-files, special protocol headers, network or physical addresses of the computer, biometric data, certificates, etc. A request to the server 200 by the client-application 220 is processed by a server application 240.\nAn authorization 260 is performed, and depending on its results, the server 200 opens a document or generates a web page, an ftp-page or other data representation based on the data type provided by the server 200. Security of data transmitted between the client 210 and the server 200 is provided by a security module 270. This conventional system can be implemented on a personal computer as well as on a web server of an AV company. Typically, the relevance of the AV databases and the effectiveness of heuristic analysis on the server side can be higher than on a personal computer.\nHowever, the principal difference in this case is the data presented by a web resource 200 being tested for malware presence based on the user identification data 250. When the AV check takes place on the server, the result of the authorization is either denial of access or a grant of the rights of the guest account. Thus, the outcome of an AV check, even at equal technological capabilities on the server and the user computer, will be significantly different.\nYet another shortcoming of conventional AV systems is restrictions imposed upon AV web-based scanners by routing rules. If a web resource is part of the network and has no external network address, it is not available outside the network, and it can only be verified by using other security tools installed on a computer system connected to this network.\nMany malware creators and hackers are aware of online scanners and knowingly block access of the scanner applications to a web resource where they have planted malware components. This makes comprehensive malware scans difficult. Furthermore, in the conventional systems, a web page cannot be scanned until it is downloaded onto a user computer. A typical web resource can contain several thousand pages. AV checking all of the web pages using a conventional system (as illustrated in FIG. 2), requires a lot of time and resources.\nThe situation is further complicated when several different web resources need to be periodically scanned. Statistics indicate that the majority of malware components and links to infected pages are located on the main (i.e., home) pages of web resources or on the first pages of sub-sections. This also complicates AV checking process, since each server often uses different identification parameters. In order to check a list of web resources, an AV application needs to have access authorization that also complicates scanning web sites for malware.\nAccordingly, there is a need in the art for a comprehensive malware scanning system that can effectively check the web resources with a minimal overhead and costs."}
{"text": "Home decor often involves the tedious application of mouldings, wallpaper, borders, and the like. Decorating a home with wall paper can become rather expensive since many rolls are required, and each roll of wall paper typically costs twenty to one hundred dollars. Furthermore, hanging the paper requires skill, patience, and time. Often, homeowners lack sufficient skill to properly hang wallpaper, or lack the time or patience to properly hang the wall paper. As a result, homeowners hire professionals to hang the wallpaper, increasing the cost to wallpaper a home. A further drawback of wall paper is that it is difficult to remove from wall surfaces when redecorating, particularly when the wallpaper is improperly hung. Yet another drawback of wall paper is that a homeowner must rely on the availability of patterns and colors, and hope that one is available which matches the style and color desired. Decorating with paint, therefore, has become an economical alternative to wall paper.\nIn part due to the reasons discussed above, painting a room has become a popular way to decorate a room. Some individuals previously considered painting as a boring option. However, now the increased availability of new colors in combination with many different methods of application can create a look quite similar to that of expensive wallpaper.\nPaint is available in a wide variety of colors. Many stores also offer mixing services, where the store employee mixes a color based on a sample which you provide. Even with these variety of colors, a person applying the paint is limited to using only one color. Alternatively, a person may apply multiple layers, creating a look containing many colors. However, this is a very time consuming approach since typically the initial layer of paint must be dry before the next layer can be applied. Alternatively, the person applying paint can utilize several different paint pans. However, having multiple pans of paint out available for use creates other disadvantages. First, significant floor space is occupied by the multiple paint pans. The person may inadvertently step into the pan and spill excess paint on shoes, clothing, and even the floor. Second, the paint in a pan not used as frequently as the others may acquire a skin on the top surface due to a drying effect. This results in impurities which remain in the pan, and eventually contaminate the roller when the paint is applied to a wall.\nOne approach to providing multiple colors of paint is taught in \"A Guide to Color & Decorating with Paint,\" published by Benjamin Moore & Co. of Toronto, Canada. A standard paint tray is provided, and a method for containing multiple colors is described. A piece of cardboard is inserted in the tray while the paint is being poured in, and the cardboard is then removed. However, this approach has several disadvantages. The paint colors may mix due to an uneven resting surface, or from agitation from the roller itself. The mixed colors create uneven results on the painted surface. Controlling the cardboard while simultaneously pouring paint is difficult. Further, the cardboard is full of paint when it is removed and is therefore an additional mess for a painter to deal with. Once the cardboard is removed, and the paint mixes due to an uneven resting surface or the pan is inadvertently kicked, the mistake of mixing the paint is irrevocable.\nAccordingly, what is needed is a paint apparatus for accommodating a plurality of colors of paint. What is further needed is a way to ensure a paint application device is properly loaded."}
{"text": "This invention relates to torque converter clutch controls and more particularly to the damping of slipping clutch controls."}
{"text": "1. Field of the Invention\nThe present invention relates to a tactile wave generating apparatus and method, and more particularly to generating amplified low frequency waves which are transmitted as tactile sound into a structure and/or to a person's anatomy. Further the present invention relates to a system where the low frequency tactile waves may be transmitted to the person's body while the full audible waves are being transmitted to the person.\n2. Discussion of the Prior Art\nElectroacoustic transducers such as loudspeakers for use in music or movie soundtrack reproduction are well known. In traditional prior art sound reproduction systems, large, powerful speakers move large amounts of air to permit a listener to feel the low frequency of sound. Listeners enjoy live concerts, in part, because they want to feel the sound pressure upon their bodies.\nIn recent years, one of the more important trends in the audio industry is that of “tactile sound” which may be described as “vibro-acoustic” or “vibro-tactile” stimulation. With tactile sound the realism of the listening experience can be enhanced by transmitting tactile waves into the person's body. For example, this could be done by vibrating the listener's seating surface of a chair or other furniture or structures. These tactile waves are able to be sensed within the person's body to add another dimension to the person's listening experience.\nThe initial applications of these devices were as sub woofer replacements or sub woofer augmentation devices. The addition of higher frequency material began to demonstrate the potential of wider bandwidth devices and the associated additional dimensions that vibro-tactile stimulation brings to the overall experience. There are many parameters that need to be evaluated when designing and/or selecting a vibro-technical device for inclusion in music and/or an entertainment system. For example, bandwidths, efficiency and power handling need to be understood. These parameters can play a big role not only in the device selection but in the amplifier selection as well.\nIt is well understood in the loud speaker industry that sufficient bandwidths (i.e., flat frequency response with sufficient low frequency and high frequency limits) is critical for high fidelity reproduction. Vibro-tactile devices, like loud speakers, are devices that must be properly designed to refine the required bandwidth for accurate response.\nIt is with these and other considerations being kept in mind that the designs of the embodiments of the present invention were created."}
{"text": "1. Field of the Invention\nThe present invention relates to a sensing unit and a method of making the sensing unit.\n2. Related Background Art\nA silicon acceleration sensor employing MEMS technology is known from the prior art as an acceleration sensor (accelerometer).\nA sensing unit of this silicon acceleration sensor is typically made through a process such as that shown in FIG. 8. More specifically, when making a sensing unit 110 of the silicon acceleration sensor, first, as shown in (a) of FIG. 8, a photomask 102 having a pattern of stripes arranged in parallel is formed on a silicon substrate 100. Deep etching (at an etching depth of 25 to 30 μm) using the photomask 102 is then performed, thereby forming a plurality of parallel trenches 104. Then, as shown in (b) of FIG. 8, an SiO2 film 106 is deposited on an upper surface of the substrate 100 and an inner wall surface of the trenches 104. Next, as shown in (c) of FIG. 8, a bottom surface part of the trenches 104 is partially removed by isotropic etching such that adjacent trenches 104 are connected to each other by the bottom surface part. As a result, a cantilever part 108 capable of reciprocating in a single parallel direction (in other words, the arrangement direction) to the surface direction of the substrate 100 is formed between the adjacent trenches 104.\nIn the sensing unit 110 having the structure shown in FIG. 8, acceleration is measured on the basis of variation in electrostatic capacity accompanying oscillation of the cantilever part (movable part) 108."}
{"text": "Although there are many personal products that one might want to have customized or made as a one-of-a kind product tailored to a particular user, a key one of these personal products is eyewear. While the invention will be described in connection with creating producing and delivering custom eyewear, it will be appreciated that the subject invention involves the creation, production and delivering of a wide variety of products that relate to the anatomical or physical characteristics of the user as well as the users preferences for particular product. That having been said, it will be appreciated that describing the invention in terms of the creation, production and delivery of eyewear carries a large number of similarities to the creation, production and delivering of a wide variety of products customized to the features and desires of the user. What follows therefore describes invention in terms of eyewear, it being understood that the invention is not so limited.\nPurchasing eyewear, while a necessity for many people, presents many challenges for consumers. For traditional in-store purchases, consumers are faced with limited in-store selection, which often requires visiting multiple stores. Yet users must explore an unmanageable array of options to find a compromise between fit, style, color, shape, price, etc. Eyewear is most commonly mass-produced, with a particular style available in one or two generic colors and sizes. Users faces are unique enough that a face can be used as a primary form of identification, yet they must choose between products made for a generic faces that are not their own. It is very difficult for users to find the one perfect pair of glasses for their unique taste, facial anatomy, and needs. They also often have difficulty visualizing what they try on because they need an optical prescription in the first place.\nRecent entrants have explored the online marketplace for eyewear in an attempt to address some of these issues. However, none of the commercially available eyewear selection systems attempt to provide a completely unique one-up, from-scratch product that is customized to the user's anatomical features, as well as the user's likes and dislikes. There is therefore a need to provide a user with a completely customizable one up product that does not rely on only off-the-shelf previously designed mass-produced or stock components. The underlying form, size, shape, or other properties of the key components must be customized to provide a truly unique and custom product for the user. Once having been able to obtain the user's image data, it is then desirable to analyze and make critical measurements of the user's face, determine user preference, and on-demand manufacture a custom piece of eyewear.\nIt is of course desirable for the process to be as automatic as possible and be one that returns to the user the most perfect one-of-a-kind piece of eyewear that he or she has ever seen. If this can be done in a relatively swift fashion, the user is provided with a quick unique piece of eyewear is that manufactured on demand.\nMore particularly, the online market is rapidly growing, though there still persist numerous problems for consumers. Consumers have poor ability to try-on glasses while shopping online. Online sites have more selection than in stores, but often the consumer is faced with endless pages of glasses from which to choose. The quality of the glasses is often unknown, and consumers are even more concerned about their new glasses fitting correctly and being comfortable since they cannot physically hold or see them until they purchase.\nA clear need exists for a shopping experience that enables a unique made-to-order product with high quality materials and design, at a price that users believe is fair and affordable for a made from scratch unique one up item, and an easier and more custom experience to creating and purchasing the perfect product for the individual, in this case a pair of glasses.\nThe concept of virtually trying on articles of clothing, including eyewear, has been discussed in the prior art for a number of years. All of the below listed patents relate to preview systems, but none relate to providing a from scratch product, relying instead on prefabricated components for a particular item.\nFor instance, Spackova in U.S. Pat. No. 4,539,585 describes a computer system to view articles of clothing on a person in an image. Mori, U.S. Pat. No. 4,730,260, and Ninomiya, et. al. U.S. Pat. No. 4,845,641, describe computer systems to virtually overlay eyewear on a person in an image. Jordan, U.S. Pat. No. 5,280,570, describes a system requiring a user to visit a store to virtually try on glasses with a realistic rendering of how their eyes will appear behind the glasses. Norton, U.S. Pat. No. 5,592,248, describes various methods of overlaying virtual images of eyewear on an image of a person's face to preview the appearance. Faye, U.S. Pat. No. 5,983,201, describes a system for users to virtually try on a variety of eyeglasses on their personal computer by connecting to an online store, selecting a subset of eyewear based on user preferences and sizes, and allowing user to purchase the frames. Gao, U.S. Pat. No. 6,095,650, describes another system for capturing an image and displaying eyewear superimposed on the user's image, including scaling of the image and detection of pupils to center the frames. Saigo, U.S. Pat. No. 6,142,628, describes another try-on system that also includes lens selection and display of lens shape in addition to frames. Waupotitsh, U.S. Pat. No. 7,016,824, describes an eyewear preview system that used a 3D face model provided by the user to overlay eyewear models on. Abitbol, U.S. Pat. No. 6,692,127, describes an eyewear try-on system that requires a wide-view camera to obtain a 3D model. Foley, U.S. Pat. No. 6,535,223, describes a system to determine pupillary distance based on an image of a person's face including an object of a known scale, as well as superimposing preview eyewear and allowing orders to be placed.\nAll of the previously described prior art explore various ways of previewing eyewear superimposed over an image of a person, but they are not on-demand systems that create, assemble and deliver a unique one-of-a-kind product from scratch. Nor do they permit previewing new custom eyewear that has not previously been mass-produced. Nor do they use the user-specific information to make eyewear better for the user. In short, they do not customize, adapt, modify, implement, or create new products such as eyewear using an on-demand system providing one-of-a-kind products from scratch. Moreover, all of the above techniques rely on previewing eyewear superimposed on the image of a person.\nOn the other hand, Fujie, U.S. Pat. No. 5,576,778, describes a system to design eyewear based on facial dimensions of a person. It is noted that Fujie is limited to controlling various anchor points on a Bezier curve that is extracted from facial image data to achieve a design. However, the specification of these anchor points or the control thereof by an individual is technical and difficult, made more so because these points are controlled using the user's words to control shape. Moreover, Fujie is limited to specifically sending polar coordinates based on Bezier curves to machine tools. This is much too complicated for a user, and the user's words alone may not be suitable as the only control.\nSoatto, U.S. Pat. No. 6,944,327, describes a system to customize eyewear based on preview images of the user's face. However, Soatto does not take into account automatically-generated user preferences. Soatto does not describe an on-demand end-to-end process and does not describe a full system that can actually manufacture eyewear. Moreover, the Soatto method is limited to specific cameras, only a frontal face image and using a method to generate a two dimensional template of the face for sizing. Limiting the preview to only a front image prevents sizing information that is critical around the temples for ensuring a good preview and comfort for the user. Moreover, most computer systems do not have multi-lens cameras conveniently available. Note, adjustment is done only through control points while maintaining a constant perimeter rim size, which is of limited application—different users will surely require different sizes. It will be appreciated that methods describing a 3D model of the face require two or more cameras not normally available to most users.\nIzumitani, U.S. Pat. No. 6,533,418, describes a system to make eyewear to order based on image previews superimposed over the user's face. However this patent only discusses changing lens shape, frame types, frame parts, and colors. It does not explain changing frame shape, but only replacing parts or changing a frame style from rimless to rimmed, which is very limiting when one wants to more fully customize eyewear. Moreover this patent does not describe automatic algorithms that size a frame to a user's face or aid in the selection of the best frames. Instead it uses a manual system like a custom order catalogue with many interchangeable parts to choose from, which could be overwhelming or too complicated for an eyewear consumer. Additionally, the preview system described only shows front and side portraits of the user with eyewear, with no interactive views, 3D views, or video, and it does not measure the dimensions of a face automatically. Further, a user is required to assist or enter information to obtain proper measurements. Finally, while the patent describes the manufacturing of eyewear, it does not clearly describe how made-to-order eyewear could actually be produced.\nWarden, U.S. Pat. No. 7,845,797, describes a method for manufacturing custom eyewear that uses a front and side image in a system with multiple cameras and lighting sources. The method requires the capture of images with and without eyewear worn on the user's face before it determines the best lens position. This method is quite limited, as it requires that the user already physically possesses the eyewear he desires, and it assumes the user simply wishes to refine the lens placement in a subsequent pair of frames. In short, this is not an on-demand end-to-end system that starts from scratch to then create, design, assemble and deliver the custom product.\nTo satisfy the needs of a typical consumer, an easy-to-use method and system that can provide a confident and enjoyable shopping experience are necessary. The system must be capable of working with the computer hardware and image capturing equipment available to typical consumer, which limits the minimum hardware to a single-lens digital camera, stand alone or embedded in a computer system, without depth or distance-measuring capability. The embodiments of this invention describe both systems to use single camera hardware and also systems that benefit from multi-camera or depth camera technology, in the event these technologies become more pervasive in a form used by consumers or in the event that a computer system is installed in a retail or office location.\nThe prior art describes technologies that are designed mostly for the aesthetic preview of the eyewear on a user. A need for a more quantitative analysis exists to enable a better experience, custom fit, custom style, automated adjustment and recommendations, and the overall ability to make an eyewear design fit with each user's unique anatomy and taste.\nOften pupillary distance is the only measurement taken to ensure the proper fit of eyewear, and that measurement alone is not sufficient to ensure a proper physical fitting of custom eyewear. More information is especially needed for advanced optics, such as progressive or digitally-compensated or freeform lenses. But regardless of the type and quantity of facial measurements needed to craft custom eyewear, the user should not be required to manually measure them. Most target users are not technologically savvy beyond following easy prompts in a web browser. A consumer needs an experience that is easier than picking and choosing parts and pieces or custom drawing every detail, especially when using only 2-D images, as the prior art has described. The method and system must enable easy customization, including automation of sizing and styles if the user desires automated recommendations. An average user should be able to obtain any eyewear design they desire and an excellent fit by having a design custom-fitted to his face, seeing a preview in a “what you see is what you get” display, and being able to make changes and see the effect on his face and fit.\nFinally, the method and system must result in a manufacturable product, such that it can be produced and sold at a reasonable cost to the user with an acceptable delivery time. It will be appreciated that a great preview system is not useful if the product being previewed is not ultimately manufacturable at a cost and in a time frame that is satisfactory to the user ordering the product.\nThus there is a compelling need for a method and system to allow greater and more personalized customization of lenses and frames, more accurate modeling and preview, more automated or assisted eyewear selection and customization, more detailed measurements, and methods to produce customized eyewear efficiently and economically to fulfill users' orders."}
{"text": "1. Field of the Invention\nThe present invention relates to a network system such as a VPN (Virtual Private Network), and more particularly to a network system that uses a firewall.\n2. Description of Related Art\nIn recent years, networks which connect multiple Intranets to each other using the internet have been proposed. The concerned networks are called extranets.\nThe VPN is known as such an extranet. By constructing a VPN, it is possible to dynamically form a work group that extends within a company or between companies. By forming a dynamic network, it is possible to link application groups in each Intranet with the internet as a base network.\nWhen constructing an extranet, flexibility and security must be ensured for each Intranet. To ensure flexibility and security, it is necessary to construct work groups between users for which access to information and services is restricted but for which there is no physical network restriction. In addition, to ensure flexibility and security, it is necessary to provide a flexible and unified interface as the interface used to provide services to work groups and to connect and disconnect users.\nAn extranet using a VPN can be viewed as a virtual extension of an Intranet. Therefore, it is preferable if the environmental conditions of each Intranet that forms a work group are the same as the environmental conditions of the Intranet operated individually. To accomplish this, it is preferable to use the same firewall as that when accessed from an outside terminal as the firewall that protects the Intranet when accessing from another internet.\nWith an Intranet operated individually, there are firewalls that use the methods shown below as firewalls for protecting the Intranet when accessing from an outside terminal. These methods can be used individually or in hybrid form to construct a firewall.\nThis method determines packets that can pass through a firewall based on IP (Internet Protocol) addresses of data subject to communication, communication port numbers that show the type of application subject to services provided to the client, etc.\nThis method executes relay or proxy response of send data for specific applications at the TCP (Transmission Control Protocol) layer level. This method uses a standard technology called SOCKS.\nThis method installs a proxy program that executes proxy responses for each protocol to operate specific applications such as HTTP (Hyper Text Transfer Protocol) or FTP (File Transfer Protocol).\nThis method registers a script for which the transactions between requests and responses have been put into rule form to a firewall for each application, and allows only packets which satisfy the rules to pass through.\nHowever, when using the methods described above as a firewall for protecting Intranets that form a work group, in other words, a firewall for protecting the Intranet from access from other Intranets, there are the following disadvantages.\nWith the packet filtering method, IP addresses are monitored as described above, but these IP addresses are set for each client (the computer that is the user terminal). Because of this, with this method, there is a fixed relationship between the client and server (computer providing services within the firewall), so it is not possible to ensure flexibility when forming a dynamic work group. Also, to perform monitoring of IP addresses for each packet, it is not possible to ensure sufficient operating capability with software installation of a firewall, making it necessary to construct a firewall with dedicated hardware.\nWith the circuit level gateway method, handling at the TCP (Transmission Control Protocol) layer is required. Therefore, when mounting a firewall on an Intranet, it becomes necessary to replace the TCP/IP library groups at the client OS (Operating System) layer. It also becomes necessary to make software changes for server applications when TCP/IP changes occur.\nWith the application gateway method, a proxy program within the firewall must be provided for each element application to construct work groups. Because of this, when adding element applications in an Intranet, it becomes necessary to newly develop proxy programs and to make firewall setting changes.\nWith the extension packet filtering method as well, as with the application gateway method, when adding element applications in an Intranet, it becomes necessary to newly register scripts and to make firewall setting changes.\nThese reasons make it difficult to use a firewall used by individually-operated-Intranets for an extranet.\nAn object of the present invention is to provide a suitable firewall for an extranet.\nTo achieve this, the network system of the present invention comprises a first Intranet and second Intranet connected to each other with the internet as a base, a first dynamic proxy server for forming a firewall to protect the first Intranet, a second dynamic proxy server for forming a firewall to protect the second Intranet, a remote access terminal connected to the first Intranet, a first object directory server installed in the first Intranet to determine if the service requested using the remote access terminal is provided in the first Intranet or provided in the second Intranet, and a second object directory server installed in the second Intranet to dynamically install a service proxy in the second dynamic proxy server when the service is provided in the second Intranet.\nWith such a structure, it is possible to protect multiple Intranets with a firewall while connecting and providing services to each other using the internet as a base."}
{"text": "A variety of guided wave optical switches have been proposed for fiber optic acoustic sensor systems. These switches include 1×1 on-off switches, 1×2 blocking switches, and 1×8 blocking switches. One characteristic requirement for all of these switches is a high extinction ratio in the off state. An extinction ratio is the ratio of two optical power levels, e.g., P 1/P 2, of a digital signal generated by an optical source, where P 1 is the optical power level generated when the light source is “on,” and P 2 is the power level generated when the light source is “off.”.\nFor different switch architectures, e.g., Mach Zehnder interferometers and delta beta reversal switches, an extinction ratio of up to 40 dB or more for each switch can be obtained with an appropriate voltage level. This extinction ratio degrades over time and temperature, resulting in voltage drift, unless the voltage level is adjusted to maintain a minimum null signal.\nOne prior art technique used to exercise an applied voltage to track and maintain the optical null requires interruption of an optical data flow in a working system. Disadvantageously, the prior art technique requires the system to be off-line.\nAnother prior art technique requires the application of a small dither voltage to the direct current (DC) port and the use of a hill climb servo to find the optimum DC bias point. Disadvantageously, this technique introduces excess noise into the system."}
{"text": "Bobrowski et al. taught that concrete remaining in the mixing drum of delivery trucks can be reclaimed for use by adding a retarding admixture to stabilize the remainder concrete (See U.S. Pat. No. 4,964,917), optionally adding new concrete to the remainder (See U.S. Pat. No. 5,203,919), and then using an accelerator just before re-using the reclaimed concrete (See U.S. Pat. No. 5,247,617).\nHines et al. taught that the amount of admixture to be dosed into the truck returning from delivery can be calculated based on remaining load size and temperature of the concrete (See U.S. Pat. No. 6,042,258) and that admixture dosing could be done on an automated basis (See U.S. Pat. No. 6,042,259).\nThe present inventors believe that prior art methods for calculating the concrete load remaining in the truck after delivery are neither sufficiently accurate nor practically convenient. For example, it is known to weigh the mixing truck on a weight scale before and after delivery (See e.g., U.S. Pat. No. 5,752,768; U.S. Pat. No. 6,123,444; U.S. Pat. No. 8,020,431; and GB 2392502), but weight can vary due to imprecision of the scale and various other factors (such as fluctuation of fuel tank and other fluid tank levels).\nIt is also known to estimate concrete discharged from the drum by counting mixing drum rotations required to discharge a known volume of concrete. A typical concrete mixing drum has a pair of mixing blades mounted on the inner drum wall, helically arrayed about the rotational axis of the drum. The blades thus function in the manner of an Archimedes' screw device. When the drum rotates in the “charge” (loading or mixing) direction, the blades push concrete towards the closed end of the drum; and, when the drum rotates in the “discharge” direction, the blades push concrete towards and through an opening located at the opposite end of the drum. The concrete expelled through the drum opening can then be guided by a chute to the desired spot where it is to be placed. Often, the load will not be fully discharged, and the remainder will be returned in the mixing drum to the plant or moved to another placement location; and the remaining volume of concrete would typically be measured by rough visual approximation or by subtracting a rough estimate of the discharged amount from the amount of the original load.\nTo this point, the current practice of estimating the remainder load has been premised upon the assumption that the amount of concrete discharged from the drum can be calculated based upon the number of drum rotations required to expel the concrete from the drum. This relation is mentioned in various patents, including U.S. Pat. No. 5,752,768 of Assh (Col. 18, line 40 et seq.), U.S. Pat. No. 8,020,431 of Cooley, and U.S. Pat. No. 8,118,473 of Compton. However, the present inventors believe this assumption is predicated on a the underlying assumption that the number of rotations required to bring concrete to the drum opening is constant from load to load, and further that the amount of discharge for each drum rotation is also constant from load to load.\nIn U.S. Pat. Nos. 6,611,755, 6,892,131, and 7,489,993, Coffee et al. disclosed that Begin Pour and End Pour events (i.e., charging and discharge) can be based upon, among other approaches, the number of drum rotations in the discharge direction, and that the truck can be weighed as part of determining amounts of concrete remaining in the drum after the End Pour event. The number of discharge revolutions for the Begin Pour event is hitherto assumed or estimated to be 1 or 2 drum revolutions, regardless of the load size of the concrete. While this may be adequate for determining Begin Pour and End Pour events, the present inventors believe that a novel system and method are required for achieving high accuracy in calculating the amount of concrete remaining in the drum after partial discharge."}
{"text": "The present invention relates to vertically disposed, rope-like structures that are ignitable from their top to dispense desired volatile materials, including, by way of example only, insect control actives such as insect repellents.\nThe fact that certain combustible materials can be ignited and will then disperse an active ingredient as they continue to smolder is well known. Incense sticks and mosquito coils are examples of products made from such materials. See e.g. U.S. Pat. Nos. 4,959,925 and 5,657,574. The disclosure of these publications and of all other publications referred to herein are incorporated by reference as if fully set forth herein.\nThe term \"fumigant\" herein means a burnable material that releases a volatile ingredient as the material burns, and preferably as it slowly smolders. A \"volatile material\" or \"volatile ingredient,\" in that context, means any material that can be released by a burning fumigant including, by way of example only, fragrances, disinfectants, and insect control actives. The term \"active\" refers to a volatile material to be released in order to achieve the desired effect of the fumigant. For a fumigant intended to repel mosquitoes, for example, an insect repellant would be an \"active.\" An \"insect control active\" is an active that repels, kills, or desirably modifies the behavior of insects. \"Insects\" herein means actual insects, as well as other small animals commonly controlled in conjunction with insects, such as spiders and the like.\nWhile conventional fumigants in the form of mosquito coils or repellent sticks are fairly inexpensive, they are often used in large quantity (e.g. thirty or so a month may be needed to control mosquitoes on an overnight basis). Further, these coils are most often used in countries having average annual income levels that are very low. In order to make such products more widely available, their cost must be reduced still further. This is particularly important in controlling the spread of malaria and certain other insect-borne diseases.\nThere have been attempts to reduce the cost of mosquito protection by burning widely available, very inexpensive, smoke-producing materials that are used without any insecticide. However, efforts to control mosquitoes by burning materials of this type, such as by burning cow dung, are not very effective and can have other undesirable characteristics, such as excessive odor.\nAcademic researchers in India have proposed the use of jute rope that has been impregnated with a mosquito repellant as a fumigant. This approach uses a relatively limp impregnated rope suspended from a hanger, which is lit on its lower end and permitted to smolder overnight. See M. Ansari et al., 31 Indian J. Malariology 57-64 (1994); M. Ansari et al., 29 Indian J. Malariology 203-210 (1992); and V. Sharma et al., 26 Indian J. Malariology 179-185 (1989).\nUnfortunately, this technique produces inconsistent results as burning rates can not be precisely controlled and, especially, can be too rapid. Also, there is some fire risk if the characteristics of the rope are not carefully controlled. For example, fire can flash up the rope, contacting structures from which the rope has been hung. Metal mesh fire guards have been shown to attempt to control this risk. Further, a portion of the rope is sometimes wasted near its upper end, as contact with the suspension device snuffs the rope.\nWax candles, having conventional fiber wicks surrounded by wax, are of course also well known. They are sometimes supported in a hollow core of a corkscrew-like wire, positioned on a flat surface. Some such candles have insect repellent (such as citronella) incorporated in the wax. However, such designs require the use of relatively expensive candle wax.\nThus, there is still a need for improved fumigants, particularly those that can dispense insecticide or insect repellent at extremely low cost."}
{"text": "The present disclosure broadly relates to the art of printing systems and, more particularly, to at least a diverter assembly, a printing system including a diverter assembly, and a method of transporting sheet media.\nKnown printing systems commonly include two or more media transport paths that divert from one another at certain points and join one another at other points. Thus, a given sheet of media can normally be transported through a known printing system along any one of a variety of transport paths.\nOne example of printing systems in which such various transport paths are utilized are those printing systems having multiple marking engines. In such printing systems, sheets of media are selectively transported from a media supply to one of two or more marking engines. Thus, a diversion point is provided along the transport pathway at which one or more sheets of media will be directed toward one of the two or more marking engines.\nUpon reaching the diversion point, a sheet of media will not itself select the appropriate media transport path along which movement of the sheet is desired. As such, mechanical diverters are typically provided immediately in front of the divergent transport paths to deflect the sheet along the desired pathway. One example of such a known mechanical diverter includes a gate that extends across the media transport path immediately in front of the diversion point of the transport path. The gate includes an upstream edge and a downstream edge, and is oriented along the transport path such that the downstream edge is pivotally supported at approximately the diversion point of the transport pathway. Thus, the gate creates a diagonally-extending blockage across the pathway that displaceable between first and second positions corresponding to sheet media diversion along the first and second transport paths.\nIn the first gate position, the leading edge of the gate is stationed away from or opposite the direction of the first transport pathway (e.g., stationed along the bottom of a horizontal transport path for diversion along an upwardly directed pathway), which thereby exposes a first side or surface of the gate. An incoming sheet of media will pass by the leading edge of the gate and contact the first side thereof, which will direct the sheet into and along the first transport path. In the second gate position, the leading edge of the gate is stationed away from or opposite the direction of the second transport path (e.g., stationed along the top of a horizontal transport path for diversion along a downward-directed pathway), which thereby exposes an opposing second side or surface of the gate. An incoming sheet of media will pass by the leading edge of the gate and contact the second side thereof, which in turn directs the sheet into and along the second transport path.\nIn operation, a printing system will transport sheets along the media transport pathway and frequently shift the gate between the first and second gate positions to selectively direct the transported sheets along an appropriate one of the first and second pathways. Commonly, a linear actuator, such as a spring-based solenoid, for example, will be operatively associated with the gate to switch the same between the first and second gate positions. One difficulty with such known arrangements, however, is associated with the continued demand for and corresponding advancement of the performance of printing systems (e.g., increased output in pages per minute). As the number of sheets transported through the pathways of a printing system increase, the number of corresponding gate switching operations is typically also increased. Thus, undesirable occurrences, such as impacts, vibrations and/or noise levels, for example, may become elevated due, at least in part, to these more frequent gate switching operations.\nAnother difficulty with known gate arrangements, which is also associated with the advancing performance of printing systems, involves the timing between the passing of a first sheet of media, the movement of the gate to a different position, and the arrival of a second sheet of media. More specifically, a given printing system will operate using a predetermined inter-document gap (IDG), which generally refers to the spacing between the trailing edge of a first sheet of media and the leading edge of a second sheet of media. However, as the output performance of printing systems continues to be improved, increasingly smaller inter-sheet gaps are expected to be used.\nIt is well known that the arrival of a second sheet of media at a diversion point prior to a gate reaching a desired gate position could result in the leading edge of the sheet of material contacting the upstream edge of the gate and thereby creating a jam or other undesirable condition. It will be recognized, them that as increasingly smaller IDGs are used, the time available for the gate to move from one position to the other is reduced. As such, the operating speed of the gate can be increased to maintain the desired operating. However, it is expected that a practical performance threshold will be eventually reached, above which only marginal increases gate switching speeds will be achievable using practical gate configurations (e.g., mechanisms of practical size or having a reasonable cost relative to the price point of the printing system).\nOne technique that can be used to increase the performance of known gate mechanisms involves initiating the switch between gate positions prior to the trailing edge of the first sheet clearing the upstream edge of the gate. However, while such techniques seem to work well at known printing system performance levels, as sheet media speeds increase and IDGs are reduced, the window for initiating the gate switch is reduced, Furthermore, care is normally exercised to ensure that that the upstream edge of the gate does not pinch or otherwise engage the first sheet of media, such as along the trailing edge thereof, and thereby undesirably slow or disrupt the movement of the sheet, This may be of particular concern where an advancement in the timing of the gate switching operations is being used to increase performance of the printing system.\nAccordingly, it is believed desirable to develop a diverter assembly, printing system and method that overcomes the foregoing and other problems and difficulties."}
{"text": "The demand for conifer trees (e.g., pines and firs) to make wood products continues to increase. One proposed solution to this problem is to identify individual trees that possess desirable characteristics, such as a rapid rate of growth, and produce numerous, genetically identical clones of the superior trees by somatic cloning. Somatic cloning is the process of producing plant embryos, in vitro, from plant cells that are not zygotes. These clones can be cultivated to yield stands, or whole forests, of conifer trees that possess the desirable characteristic(s).\nOne method for somatically cloning trees utilizes in vitro treatment of isolated, living, conifer tissue under conditions that promote formation of conifer somatic embryos, and then whole plants, from the treated tissue. The isolated conifer tissue may be cultured in the presence of one or more auxins, and/or cytokinins, to promote formation and multiplication of embryogenic tissue that is then cultured under conditions that promote formation of cotyledonary embryos that are morphologically similar to zygotic embryos produced in vivo. The embryos may then be germinated to yield conifer trees. An example of conifer embryogenic tissue are embryonal suspensor masses (ESMs) that can be formed, by tissue culture in vitro, from conifer embryos dissected from conifer seeds. By way of example, FIG. 1 shows pine embryonal suspensor masses in liquid culture. FIG. 2 shows a pine, cotyledonary, somatic embryo formed from ESM (cotyledons are visible at the top of the embryo).\nA continuing problem, however, is stimulating efficient formation of conifer, cotyledonar, somatic embryos that are capable of germinating with high frequency to yield conifer plants. Preferably, the conifer cotyledonary somatic embryos, formed in vitro, are morphologically, anatomically and biochemically similar, or identical, to zygotic conifer embryos formed, in vivo, in conifer seeds of the same species. In particular, there is a need for methods for producing, in vitro, greater numbers of zygotic-like cotyledonary somatic embryos than are produced by prior art methods. Preferably, the germination frequency and quality of the conifer, cotyledonary, somatic embryos produced by the novel methods should be higher than the germination frequency and quality of conifer cotyledonary somatic embryos produced by prior art methods."}
{"text": "The present invention relates generally to reducing open aircraft cavity acoustic resonance, and more particularly to an apparatus for expelling high frequency pulses of pressurized gas to reduce acoustic resonance within an open weapons bay of an aircraft in flight.\nThis application is somewhat related to my copending and commonly assigned patent applications xe2x80x9cJET NOISE SUPPRESSORxe2x80x9d, Ser. No. 09/973,176, xe2x80x9cHIGH FREQUENCY PULSED FUEL INJECTORxe2x80x9d, Ser. No. 09/973,161 and filed on even date herewith. The contents of these even filing date applications are hereby incorporated by reference herein.\nAs is well known, military aircraft sometimes carry weapons to be discharged during flight. The weapons are often carried within one or more cavities, commonly known as weapons bays. The weapons bays usually include a pair of doors, which are opened only when release of the weapon stores is anticipated. During the remainder of the aircraft flight, the weapons bay doors are closed, presenting a smooth surface to enhance aerodynamic performance as well as other characteristics, such as stealthiness, for example.\nWhen the weapons bay doors are opened in flight, a thin region called a shear layer is created wherein the airflow abruptly transitions from a low speed flow inside the cavity to a high speed flow outside the cavity. This shear layer is characterized by instability which causes the shear layer to form tight, circular rotating pockets of flow commonly referred to as vortices. These vortices impinge on the rear wall of the cavity causing high levels of resonance and high acoustic levels inside the weapons bay. This acoustic resonance can be strong enough to damage the aircraft or its systems, and, therefore, is quite undesirable.\nAttempts have been made in the past to reduce acoustic resonance occurring within an open aircraft weapons bay. For example, U.S. Pat. No. 5,699,981 to McGrath et al. describes a system incorporating a cylindrical member disposed on the surface of the aircraft near the leading edge of the cavity. The cylindrical member is projected into the airflow of aircraft in flight to create vortices to reduce acoustic resonance. Similarly, U.S. Pat. No. 5,340,054 to Smith et al. describes an apparatus for acoustic reduction using a series of perturbation pins placed on the surface of the aircraft for generating vortices to disrupt the shear layer, for the purpose of reducing acoustic resonances within an open cavity. While somewhat effective, these prior art systems are not without the need for improvement. For example, these patents describe systems requiring the attachment of structures onto the surface of the aircraft, increasing drag as well as significantly interrupting the desirable smooth aircraft surface. Moreover, their effectiveness decreases dramatically at aircraft speeds up to and exceeding Mach 1.\nA system recently described by Parekh, D. E. et al., Innovative Jet Flow Control: Mixing Enhancement Experiments, AIAA Paper No. 96-0308, American Institute of Aeronautics and Astronautics, AIAA, Aerospace Sciences Meeting, 34th, Reno, Nev., Jan. 15-18, 1996, includes vibrating wedges driven by piezoelectric actuators. This system, while providing the proper high frequency range of operation, is unsuited for application to the weapons bay problem because it is incapable of providing sufficient amplitude.\nAnother recent system for providing the desirable high frequency suppression operation is described in my U.S. Pat. No. 6,296,202 entitled Aircraft Weapons Bay Acoustic Suppression Apparatus, assigned to the assignee of the present invention. This acoustic suppression device utilizes an oscillatable spoiler plate which is extended into the airstream. At least one shaker is placed in operative engagement with the spoiler plate in order to oscillate the spoiler plate at high frequencies and high amplitudes to seed the shear layer with frequencies which directly complete with the natural frequency of the shear layer vortices.\nThis system represents a distinct improvement over the prior art because it provides high amplitude oscillation in addition to the high frequency operation. This desirable high amplitude operation is not obtainable with the prior art systems such as Parekh\"\"s described above.\nWhile my system represents a significant advancement over the earlier acoustic suppression systems, further improvement and refinement is desirable. More specifically, in certain situations, electromechanical devices of this type might weigh too much, or might consume too much electrical power. Also, it is well known that mechanical devices with moving parts are subject to failure and maintenance requirements. In these situations a light weight device including an injector having no moving parts might be desirable.\nA need exists therefore for an improved aircraft weapons bay acoustic suppression apparatus. Such an apparatus would provide improved high frequency acoustic resonance reduction, enhancing aircraft operation as well as aircraft longevity.\nAccordingly, it is a primary object of the present invention to provide an aircraft weapons bay high frequency acoustic suppression apparatus overcoming the limitations and disadvantages of the prior art.\nAnother object of the present invention is to provide an aircraft weapons bay high frequency acoustic suppression apparatus providing acoustic resonance reduction over a wide range of aircraft operating conditions.\nIt is yet another object of the present invention to provide an aircraft weapons bay high frequency acoustic suppression apparatus which can be utilized on a wide variety of aircraft.\nIt is still another object of the present invention to provide an aircraft weapons bay high frequency acoustic suppression apparatus that includes an injector unit for injecting high frequency pulses of pressurized gas into the airstream.\nStill another object of the present invention is to provide an aircraft weapons bay high frequency acoustic suppression apparatus including an injector unit having no moving parts.\nThese and other objects of the invention will become apparent as the description of the representative embodiments proceeds.\nIn accordance with the foregoing principles and objects of the invention, an aircraft weapons bay acoustic suppression apparatus is provided to dramatically reduce acoustic resonance within an open weapons bay of an aircraft in flight.\nAs is known in the art, opening the weapons bay doors in flight gives rise to the creation of a thin region called a shear layer where the airflow abruptly transitions from a low speed flow inside the cavity to a high speed flow outside the cavity. This shear layer is characterized by instability which causes the shear layer to form tight, circular rotating pockets of fluid flow commonly referred to as vortices. These vortices impinge on the rear wall of the cavity and correspondingly generate an acoustic wave which propagates in the opposite way, upstream. The acoustic wave thus generated interacts with the shear layer to influence the size and spacing of the vortices. If the frequency and phase of the acoustic wave coincides with that of the shear layer instabilities, the vortices generated by this xe2x80x9cforcedxe2x80x9d shear layer can become a whole number (1, 2, 3, etc.) with respect to the cavity, generating high levels of undesirable acoustic resonance. The effects of this can often be strong enough to damage the aircraft or its systems, and, as can be appreciated, are quite undesirable.\nIn accordance with the teachings of the present invention, the aircraft weapons bay high frequency acoustic suppression apparatus includes a spoiler received within the weapons bay. The spoiler is mounted such that it can be retractably extended into the airstream of the aircraft in flight. An injector unit for injecting high frequency pulses of pressurized air into the airstream is received within the spoiler. The injector unit includes a resonance tube and a nozzle having an outlet directed externally from the aircraft, into the airstream.\nResonance tubes are well known in the art. In general, resonance tubes are fluidic devices that receive a pressurized fluid input and provide a pulsating fluid output. The nature of the output, such as pulse frequency and amplitude is variable and depends on the upon the dimensions of the resonance tube.\nThe resonance tube outlet is placed in fluid communication with a nozzle. During operation of the aircraft weapons bay high frequency acoustic suppression apparatus, the spoiler is extended into the airstream of the aircraft in flight. Compressed air is directed independently into the resonance tube and into the nozzle via inlet valves. The resonance tube provides a pulsating output which is directed into the nozzle. The pulsating output of the resonance tube perturbs the flow of air in the nozzle, effectively breaking it up into discrete slugs or pulses which exit the nozzle and are directed into the airstream.\nThis high frequency discharge of high pressure gas into the airstream effectively seeds the shear layer with frequencies which directly compete with the natural frequency of the shear layer vortices even at very high aircraft speeds, something heretofore not possible with the acoustic suppression methods of the prior art.\nAdvantageously and according to an important aspect of the present invention, the injector has no moving parts, enhancing reliability as well as reducing costs. The spoiler can be made according to designs and configurations known to those skilled in this art, reducing cost as well as simplifying retrofit into existing aircraft."}
{"text": "1. Field of the Invention\nThe present invention relates to a technology for displaying information of each component that constitutes an image processing apparatus.\n2. Description of the Related Art\nIn recent years, an image processing apparatus, such as a printer, includes many functions for executing image processing, image output, and the like. Since most of such functions are complex, to see the details of such complex functions, a manual (description) attached to the image processing apparatus is required. Manuals can be broadly divided into those printed on a paper medium (paper manual) and those that are electronic in the form of a Compact Disc-Read Only Memory (CD-ROM) or World Wide Web (Web) (electronic manual).\nHowever, for the apparatus including many functions, the paper manual requires the number of pages (sheets) corresponding to descriptions of these functions. Therefore, the manual becomes huge with increased weight, and is thus difficult to carry anytime. Moreover, since the electronic manual cannot be viewed at places where a Personal Computer (PC) or the like is not installed, a user has to go to a place where a PC or the like is installed and read the electronic manual. This lacks a real-time feature.\nTo get around these problems, a technology has been disclosed in which a manual is held inside an apparatus, such as a printer, and the manual is output from that apparatus, thereby allowing manual information of the image processing apparatus to be obtained even without carrying a huge manual or installation of a PC or the like, and allowing complex functions of the apparatus to be used (see, for example, Japanese Patent Application Laid-Open NO. H8-174967). That is, in the technology disclosed this patent gazette, the user selects from an operation panel the number of a manual desired to be viewed, thereby allowing the manual to be printed out or displayed on the operation panel for viewing.\nHowever, in the technology disclosed in the gazette, the user has to search for and select a part of the huge manual where the information of the function desired to be viewed is located. This operation places a heavy burden on the user. Moreover, printing out the manual every time it is viewed increases cost. Furthermore, there is another problem in which, if the manual is displayed on the operation panel of the printer or the like, since the area of a displaying unit is small, it is difficult to see the manual.\nStill further, it will be more convenient for the user if a manual for a portion where abnormality occurs in the image forming apparatus or the like can be easily displayed on a large screen on a real-time basis."}
{"text": "1. Field of the Invention\nThe present invention relates to dynamic random access memory (DRAM) semiconductor devices and systems, and more particularly to methods and apparatus for transitioning to a self-refresh mode in a device that performs per-bank auto-refresh operations.\n2. Description of the Related Art\nDRAM devices are well known and commonly found in digital systems having a need for read/write digital memory. DRAM devices are so-named because the data in each memory cell must be refreshed periodically by reading the data, or else the stored data will be corrupted. Modern synchronous DRAM devices (SDRAMs) typically employ an “auto-refresh” mode, which refreshes one row of the DRAM memory cell array each time an auto-refresh operation is initiated by an external memory controller. An internal refresh row counter increments through the rows for successive auto-refresh operations, and wraps back to the top of the array upon reaching the bottom. The DRAM memory controller thus has some flexibility as to when it issues the auto-refresh commands to a DRAM device, as long as all rows are refreshed within the maximum time specified for the array to maintain stable data.\nMany SDRAM devices contain multiple banks of memory, with the high-order row address bits supplied to the SDRAM along with an operation determining which bank is to receive the operation. Some of these devices allow a bank address to be supplied with an auto-refresh command, and then an auto-refresh operation is performed in the bank specified by the bank address with regard to the current refresh row while a data access operation may be performed in the unselected banks at the same time. Such devices will be referred to herein as Per-Bank Refresh (PBR) SDRAM devices. The inventor of the present application has filed a copending application, U.S. Pat. application Ser. No. 11/105,169, disclosing novel PBR SDRAM architectures and methods of operation, the disclosure of which is incorporated herein by reference.\nMany SDRAM devices also incorporate a “self-refresh” mode. In self-refresh mode, the SDRAM device generally enters a lower-power state during which it does not respond to bus commands until awakened. In self-refresh mode, the SDRAM device is expected to perform its own refresh operations, based on internal timing, sufficient to retain data saved in the memory device."}
{"text": "This invention relates to the art of beauty parlor equipment and, more particularly, to a drainage tray for a shampoo bowl to facilitate washing and rinsing of a patron's hair.\nIt is of course well known that beauty parlors have shampoo bowls over which a patron's head is disposed during shampooing and rinsing of the patron's hair. Most often, the bowl has an arcuate section in a front wall thereof which cradles the patron's neck, and the rear end of the bowl includes water supply valving and a spray head for rinsing shampoo from a patron's hair following washing thereof. The shampoo bowl further includes a drain in the bottom wall thereof towards the rear of the bowl through which rinse water and shampoo drains during a rinsing operation.\nIt is not unusual for a patron to have shoulder length or even waist length hair, and the shampooing and rinsing thereof is not only difficult but potentially unsanitary to some extent. In this respect, for example, the hair of a patron with shoulder to waist length hair lies on the bottom of the shampoo bowl which makes it difficult for rinse water to penetrate the hair without the beautician having to lift the hair off of the bottom of the shampoo bowl with one hand while trying to rinse the hair with the other hand. Moreover, if the patron's hair is left on the bottom of the shampoo bowl a residue often occurs on the ends of the patron's hair, and the hair is also subject to any contaminants which might be left in the bowl from a previous patron and/or from solutions used by the beautician in connection with washing and rinsing a patron's hair. Still further, if the patron's hair is of waist length or thereabout, it can easily flow into the drain opening which not only can cause clogging of the latter but also adds to the potential unsanitary condition of the hair. Again, to avoid the latter, a beautician must support the patron's hair in one hand and try and rinse it with the other which is time consuming and difficult as well as inefficient."}
{"text": "Overlay marks are important in fabrication of semiconductor, or integrated circuit (“IC”), devices because the devices are produced by aligning several layers of conductive, semiconductive, and insulative materials one atop the other. It is critical that each layer is precisely aligned with the previous layer so that the resultant circuits are functional and reliable. If the layers are not correctly aligned, some features may be short-circuited while others may be open circuited or have an unacceptably large resistance. Typically for each technology node, an overlay error threshold is specified in the x or y direction. In other words, each layer cannot shift more than a specified distance from another layer above or below. A shift greater than the specification causes an “alignment fail”, which increases cycle time because layers may have to be reworked.\nDuring an overlay check, the position of the overlay mark on the wafer is typically sensed using a laser beam, which is bounced off the overlay mark to produce a reflected light signal reflected back to an inspector on a machine. The inspector then analyzes the reflected light signal to determine the exact position of the overlay mark. Notably, the quality of the signal reflected from the overlay mark is directly dependent on the structure and the materials. Methods to improve overlay mark signals continue to be sought."}
{"text": "LED panel lights of the above-mentioned type are known from the prior art, for which additional measures have to be taken in order to protect the layers of the layer structure and the LEDs from dust or humidity. Usually, the layer structure consists of a reflector layer, a cover glass, and a light guiding layer interposed there between which distributes the light of the LEDs, arranged on the edge of the layer structure, uniformly over the total surface. The layers of the layer structure can, for instance, be glued to one another. It is also possible to glue the LEDs to the layer structure. It is also known to use additional seals for sealing the layer structure. Also, LED panel lights of the above-mentioned type are known whose layer structure is not additionally sealed. In the course of time, these lights are subject to a deposition of dirt and/or condensed water between the individual layers so that the light output and the efficiency of the light are reduced.\nAn LED panel light is known from DE 103 15 251 B4. This light uses a sealing frame which fully surrounds an optical layer structure along the narrow sides thereof. The structure is formed of a base part, a light guide and a transparent cover. To seal the structure, the sealing frame is formed, in one embodiment, of a hot glue plastic material which permanently glues the base part and the cover together.\nIt is a drawback of the light known from DE 103 15 251 B4 that the use of a hot glue plastic material prevents an exchange of defect component parts, e.g. the LEDs. A housing frame into which the structure is inserted, and which is sealed by a second seal, is not described.\nAn illumination device for glazings in vehicles is known from US 2011/0267833 A1. In this device light-emitting diodes are used to illuminate an automotive glazing at the side. The optical layer structure is here formed of one or several transparent sheets which are glued together so as to avoid splintering. LEDs are mounted as illuminants at the side and are embraced by a seal. The seal is formed of a U-shaped metal section which is covered by rubber. Additionally, the region inside the seal, in which the LED is arranged, is filled with a transparent adhesive for sealing purposes."}
{"text": "Many monetary transactions are performed using a plastic card with a data carrying magnetic stripe. Examples are credit cards, debit cards, telephone calling cards, ATM cards and gift cards. There are other transactions (non-monetary included) that use such electromagnetically read plastic cards.\nThe problems with such plastic cards include piracy and identity theft. The British Broadcasting Corporation reports card cloning or “skimming” has doubled in the United Kingdom in the past year with the resulting thefts up to millions of dollars. See “How Credit Cards Get Cloned”, news.bbc.co.uk, Thursday, Jan. 4, 2001."}
{"text": "The present invention relates to a knife device. More particularly, the present invention relates to a multifunctional knife device which is specially used for a diving exercise.\nReferring to FIGS. 1 and 2, a conventional multi-functional knife device has a main handle 20′, a main blade 10′ disposed on the main handle 20′, an auxiliary handle 22′, an auxiliary blade 220′ disposed on the auxiliary handle 22′, and a fluorescent tube 40′ inserted in the main handle 20′. The main blade 10′ has a round aperture 30′, a back 12′, and a true edge 120′ having a plurality of serrations 11′. The auxiliary blade 220′ has a circular aperture 301′. A pivot button 31′ passes through the round aperture 30′ of the main blade 10′ and the circular aperture 301′ of the auxiliary blade 220′ to fasten the main blade 10′ and the auxiliary blade 220′ together. The auxiliary handle 22′ has a notch 23′. The main handle 20′ has a through aperture 200′ to receive an elastic element 25′, a channel 21′ to receive the fluorescent tube 40′, a groove 24′ to match the notch 23′ of the auxiliary handle 22′, and a click block 240′ inserted in the groove 24′. The auxiliary handle 22′ engages with the main handle 20′."}
{"text": "1. Field of the Invention\nThe present invention relates to a polymer film exhibiting high retardation along thickness direction (Rth); a phase difference film, a polarizing plate, and a liquid crystal display device using the polymer film; and an Rth inducing agent.\n2. Description of the Related Art\nThere have been known techniques of adding various additives to a polymer film, aiming at imparting thereto various functions.\nAn exemplary proposal has been made on a polymer film improved in ultraviolet absorption, by means of addition of a merocyanine-based compound, as a ultraviolet absorber, to the polymer film (i.e., JP-A-H08-239509). JP-A-2009-67973 discloses a polymer material containing a merocyanine-based compound as an ultraviolet absorber, a mold product manufactured by using the polymer material, and also a coated ultraviolet absorber layer. Other proposals have been made on use of a merocyanine-based compound as a wavelength dispersion modifier for polymer film (i.e., JP-A-2009-64006 and JP-A-2009-64007). JP-A-2009-270062, JP-A-2009-79213 and JP-A-2009-67983 have proposed techniques of using a merocyanine-based ultraviolet absorber combined with other types of ultraviolet absorber, aiming at improving the light stability.\nPolymer film typically used as an optically compensatory film of liquid crystal display device is sometimes required to show high Rth. Increase in the amount of addition of the additive, aimed at elevating Rth, have however occasionally degraded the stability of manufacturing, due to deposition of the additive onto the surface of the film in the process of film making. In the field of manufacturing of polymer film typically adoptable to optical films, it is therefore very beneficial to provide an additive capable of ensuring high Rth even with a less amount of addition, while ensuring a desired level of wavelength dispersion."}
{"text": "In recent years, it is possible to perform an ultrasound scan (called “an all-raster parallel simultaneous reception”) where either a plane wave or a diffuse wave is transmitted, and reflected-wave signals are received in a real-time manner by all the reception raster elements within an ultrasound frame with respect to the transmission at a time. Further, a blood flow display system (called “an ultrahigh-speed framerate method”) is known in which the all-raster parallel simultaneous reception is applied to a blood flow imaging method (called “a high framerate method”) utilizing an ultrasound scan (called “a high framerate ultrasound scan”) in which data sequences between frames are used as Doppler data sequences.\nAccording to the ultrahigh-speed framerate method, for example, when the transmission interval of the ultrasound waves at a time is 200 μs, the framerate can be calculated as 5,000 frames per second (fps). Further, Moving Target Indicator (MTI) filters are capable of processing frame data corresponding to any number of frames. In other words, because the ultrahigh-speed framerate method is capable of arranging the transmission interval of the ultrasound waves to match the framerate, it is possible to obtain an infinite time period for observation while ensuring display at a high framerate and a high folding velocity. In other words, by implementing the ultrahigh-speed framerate method, it is possible to structure a steep MTI filter having a low cut-off frequency. With these arrangements, according to the ultrahigh-speed framerate method, by eliminating as many clutters as possible, it is possible to detect blood flows in a wide range from blood flows having lower flow rates to blood flows having higher flow rates.\nFurther, as an adaptive example of the ultrahigh-speed framerate method, observing moving images at 5,000 fps in a slow playback mode makes it possible to visually recognize complex flows of blood flowing through the heart, the carotid artery, or the like. Furthermore, tracking a speckle in a blood flow makes it possible to display a flow of blood with a two-dimensional vector."}
{"text": "The field of delivering educational materials to students and professional who are not on site is referred to as “distance education” or “distance learning.” One of the concerns in this field of education is to ensure that the participant is present while the course or seminar is being presented. This concern arises because there is no proctor on site while the educational material is being delivered. The concern is exacerbated, somewhat, by the fact that more and more participants are viewing distance learning content asynchronously within their home or office without any peers to encourage attention to the presented materials. In fact, more and more content is available to stream, download or otherwise play on a single machine on a time schedule that best suits each individual participant.\nTraditionally, gauging the user's presence has been addressed by a challenge to the user to press a button or input a code into a telephone or computer to confirm presence at one or more times during or at the end of the program. While this approach ensures that the person is present, it does little to ensure that the user is paying attention during delivery of the distance learning content. As one extreme example, a course can be streamed to a computer while the user is reading and sending dozens of email messages, and, therefore, there can be great gaps in the attention given to the presented content.\nThus, while the user may be present, there have been no systems or methods developed to better ensure the user's continued attention to the distance learning content itself. Nor have any systems or methods been developed to monitor and respond to other actions taken by the participant concurrently during the presentation of distance learning content. The present invention addresses one or more of these and other problems in the art."}
{"text": "The present invention relates to a communication control unit for controlling communication between an electronic control unit and other equipment.\nSubstantial advances have been made in recent years in the field of electronic control executed by microprocessors and the like, and highly accurate control of various types of equipment has been attained. In this respect, communication between an electronic control unit which utilizes such a microprocessor and other equipment (e.g. one or more other electronic control units) is extremely important. For example in the case of recently produced types of motor vehicles, control of various units of equipment installed in the vehicle is executed by means of electronic control units. A plurality of such electronic control units can be mutually coupled by a data link so that data which is supplied to one electronic control unit from a source such as a sensor can also be supplied to other electronic control units or equipment.\nSuch communication between an electronic control unit and other units of equipment, or between on electronic control unit and another, is generally carried out by using communication control units. One communication control unit may be formed integrally with each electronic control unit, or a communication control unit may be coupled to an electronic control unit to be controlled thereby. This has the advantage of ease of implementing the communication function, while in addition, if the communication control unit is formed integrally with the electronic control unit, the additional advantage is obtained of increased general applicability. However, the following problems also arise.\nSpecifically, in the case in which a communication control unit is formed integrally with an electronic control unit, e.g., in which each electronic control unit incorporates a microprocessor and the functions of the communication control unit are implemented by the operation of the microprocessor, if a failure status occurs in the electronic control unit due to the effects of electrical noise or to occurrence of component failure, etc., then the effects of this failure status will be transferred from that electronic control unit to other equipment or other electronic control units. As a result, the trouble resulting from this failure condition may become widely disseminated to other equipment, in spite of the fact that a communication control unit is utilized in the electronic control unit which has failed.\nWhen one or more of such electronic control units are used in a motor vehicle, then such a widespread dissemination of trouble resulting from failure of a single electronic control unit can have extremely serious consequences, since various major operating function of the vehicle will generally be controlled on the basis of data produced from the electronic control units. Thus, an improvement in this respect is necessary in order to attain increased safety of operation of such a motor vehicle.\nIt is an objective of the present invention to provide a communication control unit which will overcome the problem described above, whereby in the event that a failure condition occurs in an electronic control unit, protection is provided such that this failure status will have no effect upon other equipment or other electronic control units which are connected to the faulty electronic control unit."}
{"text": "1. Field of the Invention\nThis invention relates to a spandex supply package having a cylindrical core on which is wound spandex that has a lubricating finish on its surface. More particularly, the invention concerns such a supply package wherein the amount of lubricating finish on the outer windings of spandex is considerably less than the amount of lubricating on the rest of the wound-up spandex. As a result of the decreased amount of lubricating finish on the surface of the outer windings, the stability of the package, especially with regard to sloughing-off of yarn from the package shoulders during shipping and handling, is greatly reduced.\n2. Description of the Prior Art\nSpandex is a manufactured fiber in which the fiber-forming substance is a long chain synthetic polymer comprised of at least 85% by weight of a segmented polyurethane. Most spandex is produced by dry spinning techniques and inherently is quite tacky. To avoid adherence of the spandex to itself when wound up in a supply package, as well as to avoid other problems in subsequent use of the spandex, lubricating finishes, such as silicone oils, are applied to the surface of the spandex prior to windup. Typically, the weight of the finish on the spandex amounts to in the range of 2 to 10% of the weight of the spandex.\nTo form a spandex supply package, such as a bobbin, cake, cheese, or the like, spandex is wound up on a cylindrical core. Equipment for this purpose is well known. For example, U.S. Pat. Nos. 4,398,676 (Koppen et al), 3,701,490 (Wray) and 3,409,238 (Campbell et al), among others, disclose apparatus for high speed winding of spandex thread-lines onto tubular cores to form spandex supply packages. However, during shipping and handling of such supply packages difficulties are sometimes encountered. For example, the spandex sometimes sloughs off the edges of the wound-up packages (i.e., the spandex near the flat circular faces of the wound up package falls off the edges of the package). Also, if the package is held with the cylindrical core in a vertical position during shipment, the spandex sometimes can slough off from the package and fall to the bottom of the shipping container. A similar phenomenon sometimes occurs when the spandex package is removed from the shipping container and the spandex sloughs off onto the floor. The sloughing off of spandex from a supply package creates time-consuming difficulties and waste in subsequent textile operations, such as beaming, warping, knitting, weaving, mechanical covering, air-jet entangling and the like.\nLubricating finishes for spandex and equipment for applying the finish to the spandex are well known. For example, U.S. Pat. No. 3,296,063 (Chandler) discloses certain polysiloxanes as suitable finishes for spandex. Japanese Patent Application Publication 63-66073 (Maruyama et al) describes a polyurethane elastic yarn supply package in the form of a cheese, in which the amount of finish applied to the yarn depends on the apparent elongation of the yarn within the cheese. The amount of finish generally is less on the yarns near the center of the yarn package than on the yarns near the outside of the package. Among the numerous conventional techniques used to apply the lubricating finishes to the surfaces of moving spandex thread-lines just before their windup into a supply package, are dipping, padding, spraying and the like. Known devices for applying such finishes to spandex include oiling rollers, atomizers, and the like.\nAn object of the present invention is to eliminate, or at least greatly ameliorate, the problem of spandex sloughing off the edges of supply packages."}
{"text": "Softness is a key consumer attribute of facial tissue. It is known that enhanced softness can be developed with the topical application of a polysiloxane. For nose care applications, an additional benefit to polysiloxanes can be the hydrophobicity that the polysiloxane imparts to the tissue sheet. While hydrophobicity, in general, can be an undesirable attribute for an absorbent tissue for nose care applications, such hydrophobicity can be perceived as a consumer benefit in preventing the passage of nasal secretions through the tissue and onto the users hand.\nWhile polysiloxanes can greatly enhance the softness attributes of the tissue, as well as the ability of the tissue to protect the user's hand, the ability of the tissue sheet to remove mucus and similar high viscosity materials can be reduced by application of the polysiloxane. As such, tissues treated with polysiloxane may have a reduced cleaning capability relative to an untreated tissue.\nHence, there is a need to manufacture soft tissues that have a high degree of softness and hand protection while also having the ability to effectively remove mucus from the user's nose. The effective removal of mucus from the user's nose may not only provide a cosmetic benefit in helping to clean the skin but may also provide a clinical benefit in assisting in removal of skin irritants present in the mucus. Thus, a tissue that is more soothing may also be achieved."}
{"text": "1. Field of the Invention\nThe present invention relates to a cool air circulation apparatus in a refrigerator. More specifically, the invention relates to a cool air circulation apparatus in a refrigerator having an upper curved lip formed at a front face of a lower shelf, by which the cool air contacting with ambient air is discharged outside of the refrigerator and only the cool air adjacent to a refrigerating chamber flows into the refrigerating chamber.\n2. Prior Art\nA conventional refrigerator is illustrated in FIG. 1 and FIG. 2, which comprises a freezing chamber 1 and a refrigerating chamber 2. A compressor 10 is mounted on a rear lower portion of the refrigerator, and an evaporator 20 is provided at a rear portion of the freezing chamber 2. A refrigerant is compressed by operation of the compressor 10, and the compressed refrigerant flows toward the evaporator 20, thereby cooling the circulating air by the evaporation of the refrigerant.\nFans 23, 24 for circulating cool air are provided at the rear portion of the freezing chamber 1, and the air cooled through the evaporator 20 is supplied to the freezing chamber 1 and the refrigerating chamber 2 via each cool air duct which will be illustrated later.\nThe cool air duct 25 is provided in a rear portion of the freezing chamber 1, and the cool air forcedly enters into the cool air duct 25 by the operation of the fan 23, and further enters into the freezing chamber 1 through a plurality of openings 27 formed at a duct cover 25C provided between the freezing chamber 1 and the evaporator 20.\nAnother cool air duct 26 is formed behind the cool air duct 25 opposite to the freezing chamber 1. The duct 26 is branched in two passages 26A, 26B, preferably, and each passage is further extended down along each rear side of the refrigerating chamber 2. The cool air forcedly enters into each passage 26A, 26B by the operation of another fan 24, and further enters into the refrigerating chamber 2 through a plurality of openings 28 formed at a duct cover 26C. Preferably, other openings 29 channeled from corresponding passages 26A, 26B are provided at each inner side wall of the refrigerating chamber 2.\nA cool air supply duct 42 is arranged under a partition wall 41 divided from the freezing chamber 1 and the refrigerating chamber 2, which extends from the rear portion of the refrigerating chamber 2 to the front portion of the refrigerating chamber 2. A chamber 42a for housing an air curtain fan 44 is formed at one end of the duct 42 proximal to the rear portion of the refrigerating chamber 2, and an air discharge opening 43 is formed at another end of the duct 42 opposite to the chamber 42a. The air discharge opening 43 is preferably formed along the entire width of the upper portion of an accessible opening 2C of the refrigerating chamber 2.\nThe fan 44 for generating the air curtain stream is housed in the chamber 42a, thereby enabling the air to flow smoothly. Preferably, the length of the fan 44 corresponds to the inner width of the chamber 42a, and is operated by additional motor 46.\nSince the upper surface of the duct 42 is flatly extended, and the lower surface of the duct 42 is sloped up to the air discharge opening 43, the cross-section area of the duct 42 is decreased more and more toward the opening 43. The velocity of the cool air flowing along near the upper inner surface of the air duct 42 is faster than that of the air along the lower inner surface of the air duct 42. A front end of the bent discharge opening 43 of the duct 42 is straight which causes the discharging cool air to flow straight.\nAn air collecting duct 47 is extended down a long a rear center portion of the refrigerating chamber 2, and plural air collecting openings 48, 48A, 48B channeled from the air collecting duct 47 are formed at a duct cover 47C.\nThe refrigerating chamber 2 is divided by plural shelves 49, 49A, 49B, the upper surface of which is preferably flat-shaped to enhance effective air collection. A cross-section area of the opening 48 formed at the uppermost area of the refrigerating chamber 2 is smaller than that of the opening 48A formed at the middle height area of the refrigerating chamber 2, and a cross-section area of the opening 48A is smaller than that of the opening 48B formed at the lower height area of the refrigerating chamber 2. Height of each opening 48, 48A, 48B is determined according to volume of storage foodstuffs, but each opening 48, 48A, 48B is preferably formed at approximately halfway up each shelf 48, 48A, 48B. Further, each opening 48, 48A, 48B has a rectangular shape having a long longitudinal side or an oval shape.\nThe operation of the refrigerator configurated above is illustrated as follows. When a door (not shown) is opened, the fan 44 commences operation, and simultaneously the fan 24 terminates operation. The cool air is discharged through the opening 43 by the operation of the fan 44, thus forming the cool air curtain. The cool air circulating in the refrigerating chamber 2 does not escapes from the refrigerating chamber 2, thereby maintaining a constant temperature of the refrigerating chamber 2.\nHowever, even if the cool air discharged from the opening 43 has the same low temperature between the distal layer and the proximal layer to the refrigerating chamber, the temperature of the distal layer of the air curtain increases greatly while the air stream flows downward in contact with ambient temperature air. Therefore, the air curtain flow has high temperature and the air enters into the refrigerating chamber. The high temperature air further flows into each air collecting opening, and recirculates in the refrigerating chamber. Thus, there is a problem in that the temperature of the refrigerating chamber increases, causing a decline in the cooling efficiency of the refrigerator."}
{"text": "A wireless mesh network is an autonomous collection of wireless devices that communicate with each other over wireless links. A wireless mesh network may be decentralized, and, as a result, network activity such as discovering the topology and delivering messages is performed by devices of the network."}
{"text": "The present invention relates to a sewing machines which include automatic needle threading and, more specifically, to apparatus for including in a sewing machine to provide automatic thread taking-up and threading.\nSewing machines including the capability to self-thread the needle thereof are known in the art and greatly appreciated by those using them as they are no longer required to manually insert the thread through the eye of the needle. Such a sewing machine is disclosed in Japan Published Unexamined Patent Application No. S63-89194 which is owned by the common assignee of this application. In that sewing machine, the operation of taking up the tread and threading it through the eye of the needle are simplified over manual threading. Specifically, the sewing machine semi-automatically threads the eye of the needle when an operating member provided at the front face of the sewing machine head is pressed. Since a thread taking-up member is connected to the operating member, the thread taking-up operation can be semi-automatically accomplished, too. In this semi-automated apparatus, however, the threading member, the thread taking-up member, and the operating member are interlocked using a complicated linkage mechanism and cam mechanism. The operator often has to quickly press the operating member with a rather strong force such that the linkage mechanism and the cam mechanism will operate smoothly. Unfortunately, sometimes only an experienced operator can do this successfully. In addition, the designing of the linkage mechanism and the cam mechanism and the adjusting of the interlocking of the various members is difficult and troublesome work.\nWherefore, it is an object of this invention to provide a needle-threading apparatus for a sewing machine with is fully automatic in design and operation.\nIt is another object of this invention to provide a needle-threading apparatus for a sewing machine which is simple in design.\nIt is still another object of this invention to provide a needle-threading apparatus for a sewing machine which is easy to operate and which does not require particular strength or skill on the part of an operator.\nOther objects and benefits of the invention will become apparent from the detailed description which follows hereinafter when taken in conjunction with the drawing figures which accompany it."}
{"text": "This invention relates generally to archery bow sights and, more specifically, to an archery bow sight which includes one or more sight pins securable thereto, the archery bow sight being mountable to an archery bow in a manner that allows gang movement of one or more sight pins relative to the bow.\nArchery bow sights utilizing a plurality of sight pins have been known in the art for many years. Typically, these sights utilize a bracket or other mounting structure for mounting the sight to a bow. The sight is commonly comprised of a pin plate, a pin guard, and a plurality of sight pins securable to the pin plate. The sight is mounted to a bow in a manner so that when the bow is drawn, the archer can look through a peep sight provided in the bow cable and align the tip of a pin attached to the sight with a target. For sights utilizing a plurality of sight pins having their tips vertically aligned, each individual sight pin is typically provided for aiming the bow at a target at a particular distance from the archer. For example, one pin may be positioned in the sight for aiming the bow at a target 50 yards from the archer while another pin may be positioned for a target that is at 100 yards.\nIn addition to providing individual pin adjustment relative to the sight, many prior art archery sights are mounted to the bow in a manner that allows group or what is commonly termed \"gang\" adjustment of the pins. Such adjustment is useful when the individual pins are properly positioned for a particular type or weight of arrow and shooting of a another type or weight of arrow effects the flight path. For example, for an arrow having a different quill arrangement or configuration, the pin sights may be properly positioned to account for targets at the various distances, but the flight path may be affected to the left or to the right (commonly referred to as \"windage\"). In such a situation, it may be desirable to move the entire sight to the left or the right without repositioning the individual sight pins. There are two common types of sight windage adjustment mechanisms for archery bow sights, two point and single point. For two point adjustment, the sight is typically secured to its mounting bracket at two points and may be adjusted to the left or right by loosening fasteners at these two points. With a single point adjustment sight, the sight is typically secured to a mounting bracket by a single fastener arrangement and thus can be adjusted relative to the mounting bracket by loosening of the single fastener arrangement. Similar adjustment mechanisms have been provided in the art for providing gang adjustability of the entire sight in a vertical direction.\nWhile the basic concept of securing a plurality of pins to a pin plate is known in the art, various techniques to accomplish this have been employed. For example, in U.S. Pat. No. 4,910,874 to Busch the sight pins are provided with external threads and are inserted through a longitudinally extending slot in the pin plate. An internally threaded fastener is threaded onto the sight pin in order to secure it to the pin plate. The tips of the sight pins thus extend into the sight window defined by the pin plate and the guard. By properly aligning the sight pins relative to the pin plate and the sight relative to the bow, the archer may utilize each sight pin for a target at a particular distance.\nSuch an arrangement of sight pins, however, necessarily limits the distance between adjacent tips of sight pins by the smallest transverse dimension of the structure securing the sight pin to the pin plate. For example, in the Busch patent previously discussed, the tips of the sight pins can be positioned no closer than that allowed when adjacent internally threaded fasteners are abutted against one another. For some higher power bows, it may be necessary, in order to properly sight the bow, to position the tips of adjacent sight pins closer than is physically possible if the sight pins are maintained in a parallel relationship. One approach known in the art to allow closer arrangement of the sight pins is to provide two slots in the pin plate and two different lengths of sight pins. Longer pins are secured to a back slot and shorter pins are secured to a front slot so that the tips of the longer and shorter pins extend from the front edge of the pin plate an equal distance. Accordingly, a narrow portion of the longer, back sight pins can be fitted between the wider portion of the shorter, front sight pins utilized for attaching the front sight pins to the pin plate. With such an arrangement, the distance between tips of adjacent sight pins is thus significantly decreased.\nIt is also common in the art, as shown by the sight disclosed in Busch, that a sight be comprised of several individual components including a guard, a pin plate, and various adjustment mechanisms. While each of these features provide vital features to the sight, employing separate components necessarily increases the overall weight of the sight and thus the bow to which the sight is attached. In addition, an increased number of components also increases the potential for vibration of components relative to one another which can affect the accuracy of the bow and can startle game being hunted.\nSight pins are commonly comprised of an externally threaded brass pin having a small ball formed on the tip thereof. The ball is often painted to provide for better visibility of the ball. Typically, when several sight pins are employed, each ball tip is painted with a different color indicating a different yardage for the archer. Sight pins are also known in the art which include fiber optic elements for illuminating the tip of the sight pin to make the tip more visible to the archer. The fiber optic type sight pins are commonly comprised of an elongated pin providing two transversely extending bores substantially parallel to one another, with one of the bores positioned proximate the tip of the pin. A fiber optic member having ends of a greater diameter than the bores extends from the tip of the pin, through the bore at the tip, to the second bore, and through the second bore effectively forming a partial loop on one side of the pin. The ends of the fiber optic member appear illuminated when exposed to light. With both brass and fiber optic sight pins, it is often advantageous to provide pins of different widths depending on the accuracy desired or needed.\nAccordingly, it would be advantageous to provide an archery bow sight formed from a single piece of material, integrally forming a pin guard, pin plate, and a windage or vertical adjustment mechanism. It would also be advantageous to provide a pin plate that maintains the sight pins in a substantially parallel arrangement. In addition, it would be advantageous to provide an archery bow sight that allows the sight pins to be pivoted relative to the pin plate in order to allow the tips of adjacent sight pins to be positioned relatively closely to one another. It would also be advantageous to provide a two headed sight pin that provides two sight pins of different widths."}
{"text": "1. Field of the Invention\nThe present invention relates to radio field intensity measurement devices which can display radio field intensity of radio signals. In particular, the present invention relates to radio field intensity measurement devices which change color in accordance with radio intensity of a radio signal, utilizing a chromic material, and radio field intensity detectors and amusement devices (such as game consoles, game machines, toys or the like) using the radio field intensity measurement devices.\n2. Description of the Related Art\nIn recent years, wireless devices using radio communication have been spreading due to development of electronics and the coming of an advanced information society, and have been utilized in various fields such as military affairs, medical treatment, communication, education, and commercial transaction. Radio communication using a radiowave as a communication medium takes the maximum advantages of radiowave that information can be transmitted instantly independent from time and distance, and one of the most important bases for human life and industry in modern society.\nA radiowave is a kind of electromagnetic waves, and a wave propagated in the air by energy exchange between an electric field and a magnetic field. Since a radiowave cannot be seen, there is a concern that a radiowave emitted from facilities using radiowave and wireless devices may adversely affect human bodies. Thus, there is a need to provide a device which can easily measure radio field intensity of a radiowave emitted from a radio emission device, and display information on the radiowave.\nThere are various methods for measuring radio field intensity. For example, Reference 1, Japanese Published Patent Application No. 2006-23817, discloses a radio field detector in which an antenna, a rectifier circuit and a lamp are connected and radio field intensity is detected by light intensity of a lamp. Further, Reference 2, Japanese Published Patent Application No. 2001-165973, provides an electromagnetic wave monitor in which an antennal, a storage device, and an informing means are connected, and the informing means is driven by the storage device to measure a radio field intensity.\nFIG. 38 is a block diagram illustrating a typical configuration of a radio field intensity measurement device disclosed in Reference 1.\nIn the radio field intensity measurement device 3800 illustrated in FIG. 38, a received radiowave 3820 is converted to an induction signal and the induction signal is input to a rectifier circuit 3811. The rectifier circuit 3811 rectifies the induction signal and supplied power to the lamp 3812. In other words, electricity which goes up in proportion to the intensity of the received radiowave 3820 is supplied to the lamp, and the intensity of the radiowave appears in the light intensity of the lamp which is turned on.\nIn the electromagnetic wave monitor in Reference 2, as the informing means, a light-emitting diode, a discharge lamp, or a liquid crystal display device is used."}
{"text": "Field of the Invention\nThe invention relates generally to welded aluminum-copper-lithium alloy products, and more specifically such products in the form of sections intended to produce stiffeners in aeronautical design.\nDescription of Related Art\nOngoing research is carried out to develop materials that can simultaneously reduce weight and increase the efficiency of high-performance aircraft structures. Aluminum alloys containing lithium are very beneficial in this respect, as lithium reduces the density of aluminum by 3% and increase the modulus of elasticity by 6% for each percent by weight of lithium added. In order for these alloys to be selected in aircrafts, the performance thereof must reach that of the alloys commonly used, particularly in terms of compromise between the static mechanical strength properties (yield stress, fracture strength) and damage tolerance properties (toughness, fatigue-induced crack propagation resistance), these properties being generally antinomic. Said alloys must also display a sufficient corrosion resistance, be able to be shaped using usual methods and display low residual stress so as to be able to be machined integrally.\nU.S. Pat. No. 5,032,359 describes a large family of aluminum-copper-lithium alloys wherein the addition of magnesium and silver, particularly between 0.3 and 0.5 percent by weight, makes it possible to increase mechanical strength. Said alloys are frequently referred to using the brand name “Weldalite™”.\nU.S. Pat. No. 5,198,045 describes a family of Weldalite™ alloys comprising (as a % by weight) (2.4-3.5) Cu, (1.35-1.8) Li, (0.25-0.65) Mg, (0.25-0.65) Ag-(0.08-0.25) Zr. Welded products manufactured with said alloys combine a density less than 2.64 g/cm3 and a compromise between mechanical strength and advantageous toughness.\nU.S. Pat. No. 7,229,509 describes a family of Weldalite™ comprising (as a % by weight) (2.5-5.5) Cu, (0.1-2.5) Li, (0.2-1.0) Mg, (0.2-0.8) Ag, (0.2-0.8) Mn—(up to 0.4) Zr or other refining agents such as Cr, Ti, Hf, Sc and V. Examples displayed exhibit an improved compromise between mechanical strength and toughness, but their density is higher than 2.7 g/cm3.\nPublished patent application WO2007/080267 describes a Weldalite™ alloy not containing zirconium intended for fuselage sheets (as a % by weight) (2.1-2.8) Cu, (1.1-1.7) Li, (0.2-0.6) Mg, (0.1-0.8) Ag, (0.2-0.6) Mn.\nThe patent EP1891247 describes a Weldalite™ alloy with a low alloy element content and also intended for the manufacture of fuselage sheets comprising (as a % by weight) (2.7-3.4) Cu, (0.8-1.4) Li, (0.2-0.6) Mg, (0.1-0.8) Ag and at least one element selected from Zr, Mn, Cr, Sc, Hf, Ti.\nUS Published Patent application WO2006/131627 describes an alloy intended to make fuselage plates comprising (wt. %) (2.7-3.4)Cu, (0.8-1.4) Li, (0.2-0.6) Mg, (0.1-0.8) Ag—and at least one element among Zr, Mn, Cr, Sc, Hf and Ti, wherein Cu and Li satisfy the condition Cu+5/3 Li<5,2.\nU.S. Pat. No. 5,455,003 describes a method to make aluminum-copper-lithium alloys having improved mechanical strength and toughness at cryogenic temperature. This method applies notably to an alloy comprising (in wt. %) (2.0-6.5)Cu, (0.2-2.7) Li, (0-4.0) Mg, (0-4.0) Ag, (0-3.0) Zn.\nAlloy AA2196 comprising (in wt. %) (2.5-3.3)Cu, (1.4-2.1) Li, (0.25-0.8) Mg, (0.25-0.6) Ag, (0.04-0.18) Zr and at most 0.35 Mn, is also known.\nIt was generally acknowledged in said patents or patent applications that severe homogenization, i.e. at a temperature of at least 527° C. and for a period of at least 24 hours would make it possible to achieve the optimal properties of the alloy. In some cases of alloys with low zirconium contents (EP1891247) or free from zirconium (WO2007/080267), much less severe homogenization conditions, i.e. a temperature below 510° C., were used.\nHowever, there is still a need for Al—Cu—Li alloy products having a low density and further enhanced properties, particularly in terms of compromise between mechanical strength, on one hand, and damage tolerance, particularly toughness and fatigue-induced crack propagation resistance, on the other, while having other satisfactory usage properties, particularly corrosion resistance."}
{"text": "There is a resurgence of interest in the US Navy in long range artillery as a replacement for missiles due to the cost savings involved. That is, even though the missile launcher is normally cheaper than a gun of similar caliber, each missile is much more expensive than the projectile and propellant for the gun and this soon outweighs the extra cost of the gun system.\nTo obtain greater range, the gun barrels must be longer or the chamber pressure must be greater. A combination of these two factors, together with optimum projectile design, provides the best opportunity for long range firings. However, the protrusion of the barrel above the deck is vulnerable to detection by enemy radar, hence it is of interest to reduce this to a minimum.\nTo maximize the length of barrel for a given protrusion of the barrel from the housing module, there should be a minimum of recoiled components behind the chamber/breech, and also, the recoil stroke should be as short as possible. The optimum space utilization is to have only enough room behind the back of the cannon assembly for recoiling. For this to be effected, the chamber must be separated from the barrel to allow for loading of the projectile and charge. This also allows a reduction in time between shots as the charge may be loaded at the same time as the projectile.\nThe gun described herein embodies the necessary features to allow for compact stowage as well as providing improved performance. A 5 inch caliber smooth bore barrel has been assumed, however the design can be readily scaled up or down.\nThe invention provides a gun system that incorporates a two piece cannon assembly that is separated between the chamber and the projectile forcing cone. The gun described has no traditional breech. It has a removable chamber instead which has a plug at one end that can accommodate any type of ignition system that may be required. The other end is open and is attached to the barrel with a connection fitting that is clamped to the barrel by a threaded collar. The high pressure gases created by the propellant combustion are sealed with a high-pressure gas compression seal at the interface between the chamber and the barrel. The lower side of the connection fitting has a locking collar screwed into it. The internal diameter of this locking collar has a bayonet feature or interrupted thread to suit the top end of the chamber. A torque plate, which extends out beyond the edge of the connection fitting, is attached to the lower face of the collar. A slot in the torque plate aligns with the pin of a slide actuator that is attached to a frame cross-member. When the cannon assembly is in battery position, the actuator pin is located inside the torque plate slot.\nThe connection fitting also provides the attachment for the rear cannon bearing slide which rides on rails attached to the main support structure. This also reacts any torque induced into the barrel from the actuation of the chamber/barrel locking collar.\nTo provide adequate clearance to the loading mechanism, the recoil cylinders are mounted in front of (above) the chamber/barrel connection fitting. The lower end of the cylinders are attached to the connection fitting. The structure necessary to support the front of the recoil cylinders and hence transmit the recoil force to the main support, provides a suitable mounting for the front barrel slide. The resultant extended xe2x80x9cwheelbasexe2x80x9d between the front and rear slides ensures improved shot repeatability.\nA turntable upon which are mounted two chamber elevating/support mechanisms and two projectile loader/rammer mechanisms is positioned around one of the main support structure legs."}
{"text": "I. Field of the Invention\nThe present invention relates generally to sheet metal bending brakes and, more particularly, to a sheet metal bending brake with an improved hinge.\nII. Description of Related Art\nThere are many previously known portable sheet metal bending brakes that are used in the building industry, typically for the installation of aluminum siding. These previously known sheet metal bending brakes typically comprise a frame having a planar work support surface which supports the sheet metal as well as a plurality of spaced frame members which extend over the work support surface.\nA plurality of elongated pivot arms are pivotally secured at one end to the spaced apart frame members. An elongated clamping jaw is then mounted to the opposite end of each pivot arm such that upon pivoting of the pivot arms, the clamping jaw moves toward and away from the work support surface on the frame.\nAny conventional means can be used to move the pivot arms with their attached clamping jaw between their clamping and unclamped position. In their unclamped position, the clamping jaw is spaced apart from the sheet metal support surface thus permitting the insertion and/or removal of sheet metal into the bending brake. Conversely, when the pivot arms are moved to their clamping position, the sheet metal is sandwiched in between the clamping jaw and the work support surface on the frame. Any conventional means may be used to move the clamping jaw between its clamping and unclamped positions.\nAn elongated bending arm is pivotally mounted to the frame such that the bending arm extends along the front edge of the workpiece support surface on the frame closely adjacent the front edge of the clamping jaw when in its clamped position. Consequently, with a piece of sheet metal positioned in between the frame and the clamping jaw such that a portion of the sheet metal protrudes outwardly from the front edge of the clamping jaw, pivotal movement of the bending arm in turn engages the outwardly protruding portion of the sheet metal and bends that outwardly protruding portion in the desired fashion.\nThere have been many previously known devices for pivotally mounting the bending arm to the frame. All of these previously known pivoting mechanisms, however, all suffer from one or more common problems.\nMore specifically, many of the previously known hinge mechanisms for pivotally securing the bending arm to the frame scuffed the surface of the sheet metal during the bending operation. In many situations, such scuffing is cosmetically unacceptable.\nSimilarly, many of the previously known hinge mechanisms for sheet metal bending brakes are incapable of bending the sheet metal when only a very small amount of sheet metal protrudes outwardly from the clamping jaw. For example, many previously known bending brakes are incapable of forming a bend in sheet metal of less than 3/16 of an inch in width. Similarly, many of the previously known hinge mechanisms for sheet metal bending brakes are incapable of producing hems of very small widths."}
{"text": "Document FR-2 376 919, for example, discloses a manhole which comprises a frame which is intended to be fixed in the ground and a lockable cover which closes the frame. The frame comprises an annular flat bottom plate and a substantially cylindrical skirt.\nThe cross-section of the bottom plate is substantially identical over the entire circumference thereof and, consequently, the frame is only weakly anchored in the fixing material.\nThe object of the invention is to provide a frame for a roadway device, in particular for a manhole, which allows better anchoring in the fixing material."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates to electronic systems, and more particularly to electronic systems comprising an integrated circuit transceiver.\n2. Description of the Related Art\nHand-held systems typically have a system power connection used to charge a rechargeable battery. Many newer systems support data transfer standards that specify a power connection as part of the standard. The use of a data transfer standard, such as the Universal Serial Bus (USB) Standard, in systems that are independently powered has resulted in the inclusion of at least two power pins on hand-held systems. While it is desirable that these systems be compliant with certain data transfer standards, such as the USB standard, the additional pin count is disadvantageous. Accordingly, it would be advantageous to have a method and system for reducing the number of pins needed by such a system."}
{"text": "A mobile phone is a phone that can make and receive telephone calls over a radio link, e.g., by connected to a tower of a cellular network. In addition to making calls, conventional mobile phones now support a variety of other electronic services, e.g., text messaging, email, web surfing, photography, and short-range wireless communications. Mobile phones and other mobile devices, e.g., laptops and tablet computers, are increasingly being carried everywhere by mobile device users. Since users are almost always carrying these mobile devices, the mobile devices can be used for various purposes in addition to communications purposes. For example, mobile devices that include electronic wallet application programs are gaining in popularity as an alternative to traditional wallets and physical payment cards."}
{"text": "An 802.16 family of standards have been developed by the Institute of Electrical and Electronic Engineers (IEEE) to provide for fixed, portable, and/or mobile broadband wireless access (BWA) networks (for example, the IEEE std. 802.16, published 2004 and subsequent revisions). The Worldwide Interoperability for Microwave Access (WiMAX) forum facilitates the deployment of broadband wireless networks based on the IEEE 802.16 standard. In particular, the WiMAX forum ensures the compatibility and inter-operability of broadband wireless equipment.\nIn downlink transmissions, WiMAX networks may broadcast data packets from a base station (BS) to a subscriber station (SS) or mobile station (MS), whereas in the uplink transmissions, the scheduling services may be designed to support services with different traffic characteristics and Quality of Service (QoS) requirements. A significant benefit of the converged wireless networks, such as a WiMAX network, is in the sharing of the most valuable resources—the wireless spectrum among different services. However, the wireless network convergence in a WiMAX network also comes with some challenges, due to the arbitration of uplink transmission between multiple SSs, as well as the allocation of uplink bandwidth with QoS needed for different services.\nIt is desirable to have a system and method for adaptively updating wireless uplink scheduling based on characteristics of the uplink data, to improve a user's experience, while minimizing impact on the wireless spectrum."}
{"text": "(a) Field of the Invention\nThe invention relates to a transgenic animal for the neuronal expression of HIV-1 gp160.\n(b) Description of Prior Art\nHuman immunodeficiency virus type 1 (HIV-1) is the causative agent of the acquired immune deficiency syndrome (AIDS). As a member of the lentivirinae subfamily of retroviruses, it is recognized for its ability to target the immune system and the nervous tissue (Navia, B. A. et al., 1986 b), Ann. Neurol., 19:525-535; Navia, B. A. et al., 1986 a), Ann. Neurol., 19:517-524).\nA few weeks after the primary infection with HIV-1, a burst of virus replication with a high level of viremia occurs, it is during that phase of infection that the virus likely reaches the central nervous system (CNS) (Price, R. W. et al., 1988, Science, 239:586-592). In most cases it remains clinically silent for highly variable periods of time. Occasionally, it produces an acute meningitis or meningoencephalitis. However, in a large number of AIDS patients (30-40%) the virus entry in the CNS intiates a slowly progressive dementing syndrome termed HIV-1-associated motor/cognitive complex (Janssen, R. S., 1991, Neurology, 778:773-785), which impairs cognitive and motor functions and induces behavorial disorders. At autopsy, up to 96% of these patients show neuropathological changes (Budka, H., 1990a), In: Chopra J, Jagannathan K, Sawhney IMS (eds) Advancesin Neurology. Excerpta Medica, Amsterdam, 193-202; price, R. W. et al., 1988, Science, 239:586-592; Weis, S. et al., 1993, Acta Neuropathol, 85:185-189) that typically define the HIV-1 encephalopathy (HIVE) or leukoencephalopathy (HIVL) (Villa, G. et al., 1993, J Neur Neurosurgery Psy, 56:878-884; (Weis, S. et al., 1993, Acta Neuropathol, 85:185-189).\nHIVE is mainly characterized by brain atrophy and histological changes that include white matter pallor and multiple loosely delimitated inflammatory foci disseminated in the gray and white matter (Navia, B. A. et al., 1986 a), Ann. Neurol., 19:517-524; Weis, S. et al., 1993, Acta Neuropathol, 85: 185-189). These foci are composed of microglial cells, a few lymphocytes, reactive astrocytes and very charasteristic multinucleated cells of monocyte/macrophages origin (Budka, H., 1990 a), In: Chopra J, Jagannathan K, Sawhney IMS (eds) Advances in Neurology. Excerpta Medica, Amsterdam, 193-202; Budka, H., 1990 b), Acta Neuropathol., 79:611-619; Budka, H. et al., 1991, Brain Patholo., 1:143-152). HIVL is characterized by diffuse myelin loss, gliosis and perivascular infiltration by monocytes and more rarely multinucleated microglia/macrophages. The deep white matter is preferentially and symmetrically affected (Budka, H., 1990 a), In: Chopra J, Jagannathan K, Sawhney IMS (eds) Advances in Neurology. Excerpta Medica, Amsterdam, 193-202; Budka, H., 1990 b), Acta Neuropathol., 79:611-619; Budka, H. et al., 1991, Brain Patholo., 1:143-152). In HIVE as well as in HIVL, the giant multinucleated cells are considered as the pathological hallmark and in both, morphometric studies have clearly demonstrated a significant neuronal loss (Epstein, L. G. et al., 1986, Pediatrics, 78:678-687; Ketzler, S. et al., 1990, Acta Neuropathol. (Berl), 80:92-94; Masliah, E. et al., 1992, Lab. Investigation, 66:285-291; Weis, S. et al., 1993, Acta Neuropathol, 85:185-189). In addition to these CNS changes, the spinal cord of HIV-1 infected patients may also present alterations such as myelitis or a peculiar vacuolar myelopathy for which no clear etiology has been determined yet.\nSeveral investigators have tried to identify the cells which support the expression and replication of the virus. These studies have demonstrated HIV-1 products in mononucleated and multinucleated macrophages glial cells (Cheng-Mayer, C. et al., 1987, Proc. Natl. Acad. Sci. U.S.A., 84:3526-3530; Epstein, L. G. et al., 1986, Pediatrics, 78:678-687; Gyorkey, F. et al., 1987, J. Infect. Diseases, 155:870-876; Koenig, S. et al., 1986, Science, 233:1089-1093; Stoler, M. H. et al., 1986, JAMA, 256:2360-2364; Vazeux, R. et al., 1987, American J. Pathol., 126:403-410; Wiley, C. A. et al., 1986, Proc. Natl. Acad. Sci. U.S.A., 83:7086-7093) endothelial cells and astrocytes (Epstein, L. G. et al., 1993, Ann. Neurol., 33:429-436; Wiley, C. A. et al., 1986, Proc. Natl. Acad. Sci. U.S.A., 83:7086-7093) microglial cells, but not in neurons (Savio, T. et al., 1993, J. Neurosc. Res., 34:265-272; Sweetnam, P. M. et al., 1993, Eur. J. Neurosc., 5:276-283). Furthermore, HIV-1 has been rescued from brain tissue and cerebrospinal fluid, but not from neurons (Pulliam, L. et al., 1993, AIDS Res Hum Retroviruses, 9:439-444). These findings appear rather paradoxical, considering the cognitive/motor dysfunctions, the dementing illness and the neuronal loss observed in many HIV-1 patients.\nIn absence of evidence for neuronal infectivity, several direct and indirect mechanisms have been proposed to account for the motor/cognitive disorders. a) It has been postulated that HIV-1 infection of monocytes or macrophages may activate inappropriate secretion of cytokines, such as interleukine-1 or tumor necrosis factor (TNF), that may impair neuronal and/or glial cell function or compromise the integrity of the blood-brain barrier (Wigdahl, B. et al., 1989, AIDS Res. Hum. Retroviruses, 5:369-374). b) Based on experimental observations showing that HIV-1 glycoprotein gpl20 could inhibit the growth of neurons in the presence of neuroleukin, but not in the presence of nerve growth factor (Lee, M. R. et al., 1987, Science, 237:1047-1051), it has been suggested that the HIV-1 env protein may bind to and compete for neuroleukin receptors (Brenneman, D. E. et al., 1988, Nature (London), 335:639-642; Lee, M. R. et al., 1987, Science, 237:1047-1051). As a consequence, gpl20 could directly interfere with neuronal cell function and/or cause neuronal cell death. c) Recent studies have shown that both native and recombinant gpl20, added at very low concentrations to neuronal cultures produce a striking increase in free calcium within the cells and cause cell death within 24 hours, an effect which could be abolished by adding nimodipine (100 nM), the dihydropyridine calcium channel antagonist (Dreyer, E. B. et al., 1990, Science, 248:364-367), vasoactive intestinal peptide (Brenneman, D. E. et al., 1988, Nature (London), 335:639-642), anti-gpl 20 antibodies (Dreyer, E. B. et al., 1990, Science, 248:364-367) or NMDA antagonists (Brenneman, D. E. et al., 1988, Nature (London), 335:639-642; Lipton, S., 1991, Brain Pathology, 1:193-199). The neurotoxicity of gpl20 may thus be conferred either through the NMDA receptor, via a second messenger, or directly by calcium channels (Levy, J. A. et al., 1985, Lancet, ii: 586-588; Lipton, S. A., 1992, Trends Neurosci., 15: 76-80; Sweetnam, P. M. et al., 1993, Eur. J. Neurosc., 5:276-283). Since gpl20 shares certain sequence homology with vasoactive intestinal peptide it might also compete for the same binding sites and block this neurotransmission (Lee, M. R. et al., 1987, Science, 237:1047-1051). All these studies carried out in vitro have indicated possible mechanisms by which HIV infection could lead to AIDS dementia. However, many aspects related to AIDS neurophysiopathology are still obscure and could not be approached by in vitro techniques.\nIt would be highly desirable to be provided with transgenic mice carrying the HIV-1 env gene under the neuron specific promoter of human neurofilament light gene (NFL) (Julien, J. P. et al., 1987, Genes & Development, 1:1085-1095) to further define the role of gpl20 in neurotoxicity. Such an animal model could provide some information on the effect(s) of gpl20 when expressed in neuronal cells and should help identify the mechanism(s) involved in AIDS clinical syndrome and neuropathology.\nIt would be highly desirable to be provided with transgenic animals with the neuronal expression of gpl20 and preliminary findings of the pathological evaluation."}
{"text": "1. Field of the Invention\nThis invention relates to rate-responsive pacemakers and, more particularly, to pacemakers that employ a minute volume metabolic demand sensor whose output is adjusted in accordance with a rate response dependent on the anaerobic threshold and physical condition of a patient.\n2. Description of the Prior Art\nMany attempts have been made to control the heart rate of a pacemaker patient so that it will duplicate the intrinsic heart rate of a healthy person both when the patient is at rest and when the patient is involved in various levels of exercise. Metabolic demand related parameters heretofore proposed for controlling the pacing rate include the QT interval, respiration rate, venous oxygen saturation, stroke volume, venous blood temperature, and minute volume, among others. In addition, the use of mechanical and electrical sensors which detect patient motion have also been explored in such attempts at achieving improved rate-responsiveness. Of the various parameters available, it has been found that pacemakers using minute volume as a parameter for controlling pacing rate are particularly advantageous. However, a problem with these types of pacers has been the mapping of minute ventilation to an appropriate metabolic indicated rate in a manner which accurately reflects the intrinsic heart rate.\nA further problem is that, in general, even though metabolically-related parameters used for controlling rate-responsive pacemakers, especially with the contribution of the tidal volume component of minute ventilation, react fairly rapidly to reflect changes in the patient's physical activity, the pacemakers' algorithm does not react with the same speed or time constant. This can result in the patient having a hemodynamic deficiency due to the lag time involved between the start of an increased level of exercise and the reaction thereto by the pacemaker. Of course, since linear mapping is easiest to implement, most mapping attempts so far tended to be straight line approximations, as illustrated in FIGS. 1-5. In FIGS. 1-5, both the vertical and the horizontal scales are normalized for ease of comparison. More specifically, in FIG. 1, the minimum and maximum values of the metabolic rate interval are disposed at 0 and 1 respectfully along the vertical axis. Similarly, corresponding minimum and maximum values of the minute ventilation are disposed at 0 and 1 respectively along the horizontal axis. The same convention is used in FIGS. 2-5 except that in the latter Figures the vertical axis indicates a normalized Metabolic Indicated Rate (MIR) as opposed to an interval.\nIt should be understood that rate responsive systems making use of the minute ventilation as a parameter, first calculate a long term average for the minute ventilation of a patient and then determine the difference between this long term average and an instantaneous minute ventilation obtained as described below, in conjunction with FIGS. 7-11. The resulting differential parameter is referred to as \"the minute ventilation\" for the sake of brevity. However, in the drawings, the parameter is identified as .DELTA.MV to indicate that, in fact, this parameter corresponds to the variation of the instantaneous minute ventilation from a long term average value.\nAn early attempt at mapping the minute ventilation was to use a straight linear conversion into a pacing interval, or metabolic rate interval (MRI), as shown in FIG. 1. The mapping of FIG. 1 was found to be undesirable because it led to a hyperbolic relationship between the minute ventilation .DELTA.MV and the metabolic indicated rate (MIR), as shown in FIG. 2. This relationship did not reflect the true correlation between the intrinsic heart rate and the minute volume. See \"The Range of Sensors and Algorithms used in Rate Adaptive Cardiac Pacing\" Chu-Pak Lau, PACE, Vol 15, pg. 1177, August 1992.\nAnother type of approach tried in the past was a straight linear correlation between MIR and .DELTA.MV. This method is shown in FIG. 3, and while it conformed more closely to the actual relationship between AMV and the intrinsic heart rate, it was still not satisfactory. More particularly, this straight line approach ignored important factors: first, that initially, at low levels of exercise the relationship between MV and rate is different than at a later stage after the patient has reached anaerobic threshold. Briefly, under anaerobic conditions, a patient's metabolism is operating at a level wherein the oxygen demand exceeds the oxygen intake.\nYet another mapping tried in the past was a two-slope interval-based mapping shown in FIG. 4. This method approximated the real response closely but was still not ideal.\nFinally, the mapping, referred to as augmented rate response factor, shown in FIG. 5, was also tried in order to compensate for the slow rate response at low exercise levels. In FIG. 5, line A is the rate response factor (RRF) similar to that of FIG. 3. The rate response factor is a parameter selectable in some pacemakers by a physician based on his evaluation of the patient and his own experience. Profiles B, C and D in this Figure represent three dual slope profiles, each profile consisting of two line segments: an initial, or lower rate segment for the beginning of the exercise having a first slope, and a higher rate segment having a second slope. In order to define the profiles, the physician is given the choice of selecting a low, medium or high initial RRF as the first slope. Each of these initial RRF values is higher than the base RRF. In other words, line segment B1 had a slope equal to a low initial RRF, line segment C1 had a slope equal to a medium initial RRF and line segment O1 had a slope equal to a high initial RRF. In addition to the choice of the initial RRF, the maximum and minimum MIR rates were also left to the physician's discretion or choice based on his experience and physical evaluation of the patient. Next, the break points between the line segments, indicated for each profile as B', C', and D' respectively, was arbitrarily selected at 50% of the maximum MIR, as shown. The second line segment for each profile, i.e., B2, C2, and D2 were then extended from the respective breakpoint B', C', and D' to the upper limit of the base line A, as shown. Even though improved, however, this approach still did not correspond to the actual relationship between the intrinsic heart rate and minute volume, especially at more moderate higher exercise levels. Details of this mapping are described in U.S. Pat. No. 5,292,390 incorporated herein by reference.\nIn summary, a mapping using a single RRF value as a slope is unsatisfactory because it ignores the need for a faster response during the initial or exercise level, and the need for a decreased response relative to minute ventilation at higher exercise levels. A mapping in which two line segments are used and the initial slope was set to a preselected RRF is still unsatisfactory because it relies on an artificial transition or break point between the line segments.\nAn ideal solution to the mapping problem would be to set the two slopes independently, as separate variables. However, this approach would be impractical because the transition point between the line segments would be indeterminate and, therefore, may result in responses with either abnormally extended or abnormally contracted initial slopes."}
{"text": "1. Field of the Invention\nThe present invention relates to a file folder, and more particularly to a file folder having a structure for preventing a clip from being disengaged from the file folder inadvertently.\n2. Description of the Prior Art\nTypical file folders comprise a front sheet, a back sheet, a folder strip formed between the front sheet and the back sheet and foldable relative to the front sheet and the back sheet, to receive paper sheets between the front sheet and the back sheet. A clip is further provided and engaged onto the edge portions of the front sheet and the back sheet, to clamp and retain the paper sheets between the front sheet and the back sheet.\nFor example, U.S. Pat. No. 4,904,104 to Gloeckle discloses one of the typical file folders which also comprises a clip provided to clamp and retain the paper sheets between the front sheet and the back sheet. However, a long slit or cut is required to be formed in the folder strip to receive the clip, and the clip may not be easily engaged through or into the slit or cut of the folder strip. In addition, the folder strip may interfere with the operation of the clip.\nThe applicant has also developed a file folder which is disclosed in U.S. Pat. No. 5,226,676 to Su, and which comprises three foldable strips coupled to a base sheet. One of the strips that is located close to the base sheet includes an opening formed therein to receive the clip. However, the foldable strips may be easily curled and may not be retained in position. In addition, the foldable strips may also interfere with the operation of the retaining clip.\nU.S. Pat. No. 5,285,952 to Ho discloses another typical file folder which comprises two slits perpendicular to the slit or cut as shown in Gloeckle, to slidably receive two protrusions extended from the clip, and which comprises two holes formed in the ends of each of the slits, to receive the protrusions of the clip.\nHowever, when the clip is slid away from the file folder or when the clip is completely engaged onto the file folder, the protrusions of the clip will be aligned with the holes of the file folder, and will be easily disengaged from the file folder.\nThe present invention has arisen to mitigate and/or obviate the afore-described disadvantages of the conventional file folders."}
{"text": "1. Field of the Invention\nThe present invention relates generally to vent and baffling devices for preventing the entry of rainwater and, more particularly, to a durable, rotating, water-diverting, small animal and insect entry-prohibiting sewer system ventilation cap.\n2. Description of the Related Art\nWithin the related art, numerous applications exist for devices consisting of vents with panels which prevent the entry of water. However, restrictions to the practical application of these devices is governed by the Uniformed Plumbing Code, which prohibits any such cap from restricting the egress of methane gas from the plumbing system.\nSimilarly, it is a well-known fact that a great deal of rainwater enters the sanitary sewer system during a rainfall. While a portion of this rainwater enters from improperly connected storm drain systems, some of it obviously enters via the vent system located on the roofs of homes, offices, and any structure with a sanitary sewer connection. This influx of water is then routed to the septic system or to the nearest sewage treatment plant where it must be processed along with legitimate waste water. Such large amounts of water that require treatment, especially during heavy downpours, greatly reduce the efficiency of the treatment plants and in some cases may cause the discharge of raw sewage into our lakes, rivers and streams. Over sizing of the treatment facility has been one method of dealing with this problem in the past, but this leads to inefficiencies in the plant design and causes the waste of taxpayer dollars. Additionally, in areas where developments are taking place, expansions or new treatment facilities must be added just so surges when it rains can be handled. Finally, some locations have resorted to enforcement of criminal sanctions through inspections, both physical and robotic, for the purposes of identifying, preventing, or deterring illegal dumping or other intrusion of otherwise clean rain runoff into treatments streams.\nAccordingly, there exists a need for a means by which rainwater and other undesirables can be kept from entering sanitary sewer systems by way of roof vent openings.\nIn the related art, the following patents disclose a vent and baffle unit to provide an air passage to an attic. These include U.S. Pat. No. 4,269,007 issued in the name of Ward, U.S. Pat. No. 4,214,510 issued in the name of Ward, and U.S. Pat. No. 4,125,971 issued in the name of Ward.\nU.S. Pat. No. 5,830,059 issued in the name of Sells describes a ventilating cap for a roof including a cover, a baffle, batts, or foam rubber blocks.\nU.S. Pat. No. 5,655,964 issued in the name of Rheault discloses a static ventilator with louvers, a triangular base, and a cap for preventing anything from falling into it.\nU.S. Pat. No. 5,549,513 issued in the name of Thomas et al. describes a roof structure with a ventilation device and a plurality of baffles for preventing the entry of rainwater.\nU.S. Pat. No. 5,421,776 issued in the name of Sakamoto et al. discloses an exhaust air hood for reducing backflows.\nU.S. Pat. No. 5,201,879 issued in the name of Steele describes a vent for enclosures that prevents the entry of water.\nU.S. Pat. No. 4,899,647 issued in the name of Garries et al. discloses a ventilator system including an exterior panel and an interior baffle to prevent water from draining into it.\nWhile ventilating cap and static ventilator devices are incorporated into this invention in combination, other elements are different enough as to make the combination distinguished over this related art.\nConsequently, a need has therefore been felt for an improved but less complex mechanism that keeps rainwater out of sanitary sewer systems thereby reducing the loading on water treatment facilities and reducing pollution.\nIt is therefore an object of the present invention to provide an improved sewer system ventilation cap.\nIt is a feature of the present invention to prevent the entry of rainwater into sanitary sewer systems thereby reducing the loading on water treatment facilities.\nIt is another feature of the present invention to prevent undesirables such as small animals, birds, insects and the like from entering the system.\nBriefly described according to the preferred embodiment of the present invention, a durable, rotating, water-diverting, small animal and insect entry-prohibiting sewer system ventilation cap is disclosed, comprising a hollow main chamber forming a 90xc2x0 angle having an anterior end, a posterior end, and an upper portion. The upper portion of the main chamber is designed such that it measures 45xc2x0. The 45xc2x0 angle design prevents the entry of water through the main chamber. A screen of wire mesh is disposed within the interior of the anterior end of the main chamber to prevent the penetration of physical debris or particulates into the sewer system ventilation cap.\nThe posterior end is designed so as to slidably receive a secondary chamber of a generally linearly elongated tubular configuration. The secondary chamber comprises a first end and a second end. The first end slidably engages the posterior end of the main chamber in such a way to allow the main chamber to rotate 360xc2x0 around the secondary chamber when confronted by wind or storms. The main chamber\"\"s ability to rotate 360xc2x0 allows for the evacuation of sewer gases, allows proper ventilation and prevents the entry of rain water through the main chamber. The warmth of the sewer gas prevents the main chamber from freezing.\nThe lower end of the secondary chamber is connected to a plumbing vent of a conventional sewer system.\nIn accordance with a preferred embodiment, a ventilation cap having grill work with a cover over the grill work is envisioned. The cover is of a dome-like configuration having a diameter so as to effectively shield the plumbing vent from the entry of rainwater and simultaneously allowing the sewer gases to vent. The grill work is comprised of several grill members of a linearly elongated rectangular configuration having a front end and a back end. The front ends of the grill members are coupled to the underneath portion of the cover and are collectively aligned so as to form a circular configuration. The back ends of the grill members are attached to the plumbing vent of a conventional sewer system.\nIn another preferred embodiment of the present invention, a sewer system ventilation cap having baffles, drain holes, and an interior lip is envisioned. The baffles are of an elongated oval shape defining a hole therethrough and are located along the external circumferential surface of the sewer system ventilation cap. The baffles allow for the evacuation of sewer gases and allows the system to vent properly. The interior lip is located at the base and inside of the sewer system ventilation cap. Drain holes rest along the interior lip so as to allow the exit of any rainwater that has entered.\nAn advantage of the present invention is that it conserves natural resources and reduces pollution.\nAnother advantage of the present invention is the savings in time and manhours represented by the decreased water throughput required by a wastewater treatment systems that can be achieved if a significant number of structures incorporate the teachings of the present invention.\nYet another advantage of the present invention is that it postpones the expansion of existing sewer treatment facilities or prevents the construction of new sewer treatment facilities, in turn saving tax dollars."}
{"text": "Vertically adjustable sights, for example, those used in the field of archery, are known to be adjustable to account for many external factors, e.g. the distance to the target. Some current sights use cam members or other such mechanisms to adjust the sighting element, such as a sight pin, which can be inaccurate and/or hinder the ability to lock the sight pin at a select adjustment. Some current sights are also not adaptable to many different sizes of firearms, bows or other similar type of weapon or equipment.\nOne sight designed to overcome such problems is disclosed in U.S. Pat. No. 7,360,313, which is hereby incorporated by reference herein in its entirety for all purposes, and shares common inventorship and ownership with the present invention. The sight disclosed in U.S. Pat. No. 7,360,313 utilizes a gearing system to adjust the sight pin, includes a locking mechanism to prevent the sight pin from undesired movement and is adaptable for number of different sized bows. However, the sight includes a number of moving pieces to translate the desired adjustment as indicated by the use of an adjustment mechanism, e.g. a dial, to the sight pin.\nHaving many moving parts in such a sight increases the cost to manufacture both from a materials and assembly perspective. Having many moving parts also increases: (1) the failure rate by having an increased number of parts that could have defects, (2) the wear on the moving parts and (3) the opportunity for external elements to affect a part of the mechanism, e.g. dirt or debris. For example, having a number of slots within which selected parts move presents multiple opportunities for dirt, debris, rain, snow or other elements to interfere with the operation thereof. Having many moving parts can also increase the weight of the sight and, thereby, the weight of the object to which the sight is attached, for example a bow, which can affect accuracy due to fatigue in holding the bow.\nAs such, there is a need for a selectively lockable, geared, adjustable sight that has a minimum number of moving parts.\nIt will be understood by those skilled in the art that one or more aspects of this invention can meet certain objectives, while one or more other aspects can lead to certain other objectives. Other objects, features, benefits and advantages of the present invention will be apparent in this summary and descriptions of the disclosed embodiment, and will be readily apparent to those skilled in the art. Such objects, features, benefits and advantages will be apparent from the above as taken in conjunction with the accompanying figures and all reasonable inferences to be drawn therefrom."}
{"text": "The present disclosure relates generally to the battery systems for vehicles deriving at least a portion of their motive power from an electrical power source. More specifically, the present disclosure relates to systems and methods used to pre-charge a bus capacitance in vehicles using parallel battery packs.\nThis section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present disclosure, which are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present disclosure. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nVehicles, such as cars, trucks, and vans, are widely used to facilitate the movement of people and goods in modern society. Vehicles may utilize a number of different energy sources (e.g., a hydrocarbon fuel, a battery system, a capacitance system, a compressed air system) to produce motive power. For example, certain vehicles may include a battery system and use electric power for part or all of their motive power.\nIn general, such vehicles may provide a number of advantages as compared to traditional, gas-powered vehicles that solely rely on internal combustion engines for motive power. For example, such vehicles may produce fewer undesirable emission products and may exhibit greater fuel efficiency as compared to vehicles using only internal combustion engines to propel the vehicle.\nTo conserve stored energy, the battery systems used to provide motive power for vehicles may be disconnected from the vehicle's power systems when the vehicle is not in operation. Because of the large voltage difference and the bus capacitance between the battery systems and the disconnected power systems, immediately reconnecting a battery system to a disconnected power system on the vehicle may cause a large inrush of current into the bus capacitance, potentially damaging both the battery and the power system. Hence, pre-charging systems have been developed that are designed to limit the inrush current of a battery system until the bus capacitance is sufficiently charged. These pre-charging systems often include a pre-charge resistor that may be connected in series between the battery system and the bus capacitance until the bus capacitance is charged. The capacitance of the bus capacitance and the resistance of the pre-charge resistor effectively form an RC circuit, causing a delay in the pre-charge system in relation to the RC time constant of the RC circuit."}
{"text": "The present invention relates to cooling mechanisms and, more particularly, to an apparatus for containing and cooling an integrated circuit device.\nIt is well known that the performance of some electronic devices can be improved by cooling them to a low temperature, and even to a temperature below 0.degree. C. As a result, a variety of devices have been designed to cool integrated circuit devices. For example, U.S. Pat. No. 4,730,665 discloses an apparatus for cooling a high-density integrated circuit device. This apparatus includes a housing in which the circuit device is sealingly mounted in an internal chamber of the housing. The chamber is at least partially filled with a plurality of thermally conductive spheroids which are biased into thermally conductive contact with each other and with the integrated circuit device to maximize the heat exchange surface. This apparatus further includes an inlet and an outlet for passing a dielectric immersion coolant through the chamber in direct heat exchange relationship with the spheroids and with the high-density integrated circuit device. Another example is shown in U.S. Pat. No. 4,758,926 which describes a device for enclosing, protecting and cooling semiconductor integrated circuit chips. The device includes a generally planar substrate with the chips positioned thereon. Signal connections are provided between at least some of the chips. A heat sink is positioned in contact with the chips and includes microchannels through which a cooling fluid flows for purposes of transferring heat generated by the chips to such fluid. Manifolds are provided to direct the fluid to and from the microchannels, and microcapillary slots may be formed on the heat sink surface adjacent the chips to receive liquid to generate attractive forces between the heat sink and chips to facilitate heat transfer. Circuitry is provided to distribute power through the device and to the chips. Yet another example is shown in U.S. Pat. No. 5,032,897 which discloses a thermoelectrically cooled integrated circuit package which includes an insulative module which defines a cavity. A thermoelectric cooler is positioned within the cavity, and an integrated circuit chip is connected to the thermoelectric cooler.\nHowever, a problem which may be encountered with the aforementioned devices is the formation of frost on the surface thereof during low temperature operation, Such problem may result in water developing on the printed circuit board which may cause structural damage thereto or shorting of the electrical circuitry thereon. Also, the aforementioned devices may be relatively difficult and expensive to manufacture."}
{"text": "Thrombocytosis is a condition in which there are an increased number of platelets in the peripheral blood. Thrombotic thrombocytosis, in which the increased platelet count results in intravascular clotting, is most commonly associated with chronic myeloproliferative diseases. Examples of myeloproliferative diseases include primary thrombocythemia, polycythemia vera and chronic myelogenous leukemia. Primary thrombocythemia, also known as essential thrombocythemia, idiopathic thrombocythemia, hemorrhagic thrombocythemia and megakaryocytic leukemia, may result in spontaneous hemorrhages, both external and into the tissues. Polycythemia vera involves both an increase in red cell mass and an increase in total blood volume. An increase in circulating platelets is also often associated with reactive processes. In reactive processes to splenectomy, surgery, inflammation, iron deficiency or hemorrhage, the platelets typically exceed one million per microliter. This level of thrombocytosis is accompanied by little tendency toward thrombosis unless there are additional risk factors such as preexisting arterial disease or prolonged immobility. In chronic myeloproliferative disease, circulating platelets often approach 5 million per microliter. In addition to hemorrhage, thrombosis is often observed in association with such a high number of platelets. It has been postulated that thrombosis is due to spontaneous platelet aggregation in the circulation of patients with chronic myeloproliferative disease. Left untreated, the thrombosis may lead to ischemic lesions of the toes, pulmonary embolism, apoplexy, and even death.\nCurrent clinical practice in these neoplastic-thrombocytosis patients is to lower the platelet count when values exceed one million per microliter. The treatment of choice is presently concomitant platelet pheresis in association with marrow suppressive therapy using radioactive phosphorus or alkylating agents. The disadvantages to this concomitant therapeutic approach is that the former therapy requires hospitalization and technical assistance and the latter treatment is well recognized to enhance leukemogenic responses.\nA more desirable pharmaceutical approach to this condition would be administration of a drug with a direct, specific effect on lowering the platelet population without any side effects. Currently, at least two drugs with the desired thrombocytopenic effect are being evaluated in the treatment of thrombocytosis associated with chronic myeloproliferative disease: hydroxyurea and anagrelide. In the case of hydroxyurea, granulocytopenia is an undesired side effect and there is also a theoretical concern that this drug may have leukemogenic potential. Anagrelide, 6,7-dichloro-1,5-dihydroimidazo-[2,1-b]quinazolin-2(3H)-one monohydrochloride, which is currently being given 4 times a day with maintenance doses of 1.5 to 4.0 mg/day, has been reasonably effective in lowering platelet counts beginning about 5 days after initiation of treatment. However, in addition to transient side effects of headache and nausea, higher doses of this drug may cause potent aggregation. U.S. Pat. No. 4,743,445, indicates that alpha interferon may be used to treat thrombocythemia, but that it also reduces the white blood cell count."}
{"text": "The present invention relates to a lock assembly, and more particularly to a mounting arrangement for a universal core assembly into multiple lock housings.\nNumerous types of conventional lock assemblies are utilized for various applications. Homes and commercial establishments are protected predominantly by key-actuated pin tumbler locks. In a typical lock, a core assembly houses a rotational cylindrical plug having a longitudinally extending keyway. A driving member such as a cam is connected to the rear face of the plug. Rotation of the plug rotates the cam, which thereby rotates a driving member. The driving member actuates a bolt-throwing or latch-moving mechanism.\nThe interface between the plug and the case is called the shear line. A plurality of radially extending, parallel chambers is formed in the case and the plug. Spring-biased pins are disposed in each chamber. Under normal conditions, the drivers block the shear line, thereby preventing the plug from being rotated relative to the case. However, when a properly configured key is inserted into the keyway, the drivers and lower pins are moved so that the top of the lower pins and the bottom of the drivers meet at the shear line. The plug can then be rotated to cause rotation of the driving member and subsequent retraction or extension of the bolt or latch.\nLocksmiths frequently must re-key or replace residential or commercial locks on short notice. To this end, interchangeable core assemblies are manufactured by various lock makers. Disadvantageously, the interchangeable core assemblies are relatively complicated. Typically, the interchangeable core assembly, even those from a single manufacture, is specific to a particular lock type and including mounting structure specific thereto. For example, a knob lock assembly, a lever lock assembly, and deadbolt lock assembly each utilize a core assembly and mounting arrangement particular to a knob, a lever, and a deadbolt, respectively. Such an arrangement complicates re-keying and replacement of residential and commercial locks.\nAccordingly, it is desirable to provide an uncomplicated mounting arrangement for a core assembly that is readily mounted into multiple lock types."}
{"text": "Healthcare providers, such as doctors and nurses, frequently use handheld devices when providing healthcare. Many of these handheld devices include electrical devices that must be powered by electricity. One example of such a handheld device is an otoscope, which includes a light to illuminate a patient's ear canal during an examination.\nSome handheld instruments include a power cable that can be plugged into a wall receptacle to deliver power to the handheld device. The power cable limits the use of the handheld instrument to locations near a wall receptacle. The power cable is also inconvenient because it can inhibit free movement of the handheld device.\nBatteries are sometimes used in handheld devices to overcome the drawbacks of power cables. While batteries do allow the healthcare provider to freely move the device independent of a wall receptacle, batteries have their own limitations. Batteries must be periodically recharged or replaced. Batteries are also slow to recharge, and so the handheld device may be out of service for some time. In addition, batteries have a limited life, and can be harmful to the environment if not disposed of or recycled properly."}
{"text": "Automatic card shuffling machines were first introduced by casinos approximately ten years ago. Since then, the machines have, for all intents and purposes, replaced manual card shuffling. To date, most automatic shuffling machines have been adapted to shuffle one or more decks of standard playing cards for use in the game of blackjack. However, as the popularity of legalized gambling has increased, so too has the demand for new table games utilizing standard playing cards. As a result, automatic shuffling machines have been designed to now automatically “deal” hands of cards once the cards have been sufficiently rearranged.\nFor example, U.S. Pat. No. 5,275,411 (“the '411 patent”) to Breeding and assigned to Shuffle Master, Inc., describes an automatic shuffling and dealing machine. The '411 patent describes an automatic method of interleaving cards as traditionally done in a manual fashion. Once interleaved, the entire stack of shuffled cards is positioned above a roller that removes and expels a predetermined number of cards from the bottom of the stack to a card shoe. Once the predetermined number of expelled cards are removed from the shoe by a dealer, a second set of cards is removed and expelled. This is repeated until the dealer has dealt each player his or her cards and has instructed (e.g., pressed a button on the shuffler) the shuffling machine to expel the remaining cards of the stack.\nThe '411 patent and related shufflers, having a dealing means, suffer from the same shortcomings—slowness, misdeals and failure. However, the machines currently marketed are still favored over manual card shuffling. On the other hand, since casino revenue is directly proportional to the number of plays of each wagering game on its floor, casinos desire and, in fact, demand that automatic card shufflers work quickly, reliably and efficiently.\nAccordingly, the present invention utilizes a proprietary random card ejection technique, in combination with a novel card separation and delivery unit, to overcome the aforementioned shortcomings. The present invention uses random ejection technology to dispense individual cards from a card input unit to a card separation and delivery unit of the shuffler. A card stop arm and floating gate control the number of ejected cards that may, at any one time, travel to the card separation and delivery unit. The ejected cards are then separated by a feed roller system which propels the cards to a collection unit. Once a predetermined number of cards are propelled to the collection unit, additional cards are ejected from the card input unit. A shuffler processing unit in communication with internal sensors controls the operation of the shuffler.\nAn audio system is adapted to communicate internal shuffler problems and shuffler instructions to an operator. Preferably, the audio system is controlled by the shuffler processing unit in communication with a second local processing unit."}
{"text": "Generally, in language, a word is the smallest element with meaning. A written language is the representation of a language by means of a writing system. An alphanumeric writing system may use a set of symbols, letters, and/or numbers, to form a word. In a logographic writing system, a logogram, which is a single written character, is used to represent a complete grammatical word or morpheme. A morpheme is the smallest semantically meaningful unit in a language. For example, some Chinese characters are logograms.\nText is the representation of written language. Printed text can be scanned, for example, using optical character recognition (OCR), to create an electronic image of the text. OCR is the electronic conversion of scanned images into machine-encoded text. The converted machine-encoded text may then be electronically searched and/or used in various machine processes, such as text mining, machine translation, etc. When running an OCR application on a scanned image, boundary information for the text is created. In character recognition, boundaries can be a real or imaginary rectangle which serves as the delimiter between consecutive letters, numbers, and/or symbols in alphanumeric words and between lines in character words (e.g., Chinese character words). The boundary information can include the rectangular coordinates for the lines that make up Chinese character words, and letters, numbers, and/or symbols in alphanumeric words.\nTypically, when a scanned image is of poor quality or if the scanned image contains logographic characters (e.g., Chinese characters), the OCR application may make mistakes in detecting the boundaries, and applications and processes, which may rely on the boundary information, may generate incorrect results. For example, a Chinese character word or an alphanumeric word may be split into multiple parts, causing typographical and grammatical errors. Editors may spend a significant amount of time in trying to detect and correct the errors."}
{"text": "A user of a video camera may be moving, have unsteady hands or hold the camera at one or more different angles while shooting a video. As a result, a recorded video may include one or more frames of the same scene that are rotated at different angles with respect to a reference (e.g., horizontal or vertical). The rotated frames may make the recorded video difficult to view or edit."}
{"text": "Today, some software programs contain one million lines of code (LOC) or more. In order to ensure the reliability of these software programs, the software programs undergo testing before releasing the software programs to consumers. A first phase of testing is commonly referred to as verification testing. Verification testing attempts to find errors in the source code of the program. The identified errors may then be corrected before the software program undergoes another phase of testing, commonly referred to as run-time, or beta testing.\nOne approach for performing verification testing of software programs is based on dataflow analysis. In dataflow analysis, a property that is being tested is analyzed through each path of the program. The analysis determines a state for the property at each point in the program and maintains information regarding the state of the property. Upon traversing two paths of a branch statement, the analysis may determine that the property could be in multiple states at the end of the branch statement. These multiple states are combined and carried forward in the analysis. Later, when another branch statement occurs, the analysis traverses both paths of the branch statement using the combined multiple states. This method of analysis is extremely efficient. However, the results are not very accurate. Typically, the results indicate an error in the source code even when the source code does not have errors. A programmer is then responsible for reviewing the source code and for determining which reported errors are actual errors. Even though dataflow analysis provides inaccurate results, it is commonly performed on both small and large software programs.\nAnother approach for performing verification testing of software programs is commonly referred to as path-sensitive analysis. Briefly, path-sensitive analysis tracks a concrete state of the program (e.g., values of variables). Therefore, when a branch statement occurs in the program, path-sensitive analysis uses the concrete state to determine which paths of the branch statement are feasible. The feasible paths are then traversed. Whenever the analysis cannot determine which path of the branch statement to traverse based on the concrete state, the analysis duplicates the concrete state so that both paths may modify their copy of the concrete state as needed. This duplication of the concrete state at each branch statement greatly increases the search space and is commonly referred to as “search space explosion”. For example, if a program has one hundred (100) branches, the search space may grow exponentially by 2100. This exponential growth of the search space limits path-sensitive analysis to small software programs.\nUntil the present invention, full-scale reliable verification of a software program having a large code base was unattainable."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of secure software rental systems.\n2. Background Art\nThe use of personal computers (PC) and communication modems (CM) has increased significantly in recent years and is expected to continue to grow. Using electronic mail capabilities, PC users communicate more frequently amongst themselves. Additionally, centralized computer systems have evolved allowing PC users to access large databases. Such databases include various information libraries: news, weather, sports, stock markets, entertainment, education, and so on. Access to such databases is commonly controlled so that users must subscribe to the centralized computer systems. In a typical session, the user connects to the centralized computer system using the PC, transfers information to the user's PC, and is further useable without being connected to the database of the centralized computer system. The centralized computer system enables a large number of users to concurrently access the database of the central computer system.\nWhile centralized computer systems frequently provide access to information databases, such systems less frequently provide access to copyrighted application software. The primary reason for not providing copyrighted application software from databases of centralized computer systems is due to a lack of tamper-proof security methods and apparatuses for preventing unauthorized copying of copyrighted application software. Prior art systems do not provide a comprehensive method or apparatus for permitting the rental of copyrighted application software without having any possibility of the copyrighted application software being copied and used without being connected to the database.\nA prior art system, disclosed in U.S. Pat. Nos. 4,796,181 and 5,047,928 issued to John D. Wiedemer on Jan. 3, 1989 and Sep. 10, 1991, respectively, implements a computer software security and billing system that enciphers an application program using a numeric key. The computer of the user requires a hardware security device and a removable billing device. Both devices carry unique codes. The security device containing the billing device is coupled to the user's computer. A security program accesses the application software and writes billing information into the billing device. The billing module must be periodically replaced so the user can be charged for the software usage. Thus, the system of Wiedemer is directed to a security device including a billing device that is installed in a user's computer for enciphering/deciphering software and billing for usage of the software. This system disadvantageously requires special hardware for billing use of application software and does not use a dynamic password for preventing unauthorized use of application software.\nAnother prior art system, disclosed in U.S. Pat. No. 4,999,806 issued to Fred Chernow, et al., on Mar. 12, 1991, is a system for distributing software by telephone. A central station accepts credit card information, transmits an acceptance code to a caller, and terminates the call. The central station first verifies the caller's credit card, and then calls back the caller. The transaction is continued after receiving the acceptance code. The central station transfers a control transfer program and initialization program to the caller. The caller (or purchaser) executes the initialization program so that the central station can control the caller's computer. The control transfer program then transfers a protection program for ensuring that a copying program is not resident in the memory of the caller's computer. A storing program is then transferred to the caller's computer for modifying the purchased program for storage on the caller's computer. The purchased program is then transferred to the caller's computer. During execution of the system for distributing software, the various transmitted programs are erased so that only a copy of the purchased software remains on the caller's computer. Thus, the system of Chernow, et at., is directed to a system of transmitting copy protected versions of software to a caller's computer for a limited amount of time similar to a demonstration. The system of Chernow et al., is similar to copy protection of software and does not use a dynamic password for preventing unauthorized use of application software.\nA further prior art system, disclosed in U.S. Pat. No. 5,138,712 issued to John R. Corbin on Aug. 11, 1992, implements a method and apparatus for licensing software on a computer network. Encrypted license information is stored in a license token, and is sorted in a database controlled by a license server. To access a program, the license server locates the correct license token for a software application and transmits the license token to a license library. The application has an attached application specific license access module that decodes the licensing token. The license information is verified by license library routines coupled to the software application. The license is then checked out and the license token is updated. The application specific license access module encodes the updated license token before returning it to the license server. Thus, only a single application can be breached by unauthorized cracking of an encrypted application. Thus, the system of Corbin is directed to providing network protection against unauthorized use of software in a computer network.\nThus, the prior art is not able to provide a dynamic secure software rental system."}
{"text": "This invention relates to a method and system for the preparation of specimens for analysis by electron microscopy, and more particularly to a method and system for the milling of specimens using an ion beam.\nIon beam milling systems are used for the preparation of specimens whose internal and surface structures are then analyzed using transmission electron microscopy (TEM) or scanning electron microscopy (SEM). Many techniques have been developed which have included the use of mechanical cutting, grinding, and/or polishing operations in combination with ion beam bombardment. In many instances, such techniques have required that the specimen be repeatedly moved from one apparatus to another (either different preparation devices or observation devices) and/or be moved or realigned while mounted in a cuffing, grinding, or milling device.\nTransmission electron microscopy is an important technique for studying the detailed microstructure of many materials. Improvements in the design and operation of electron microscopes have generated considerable interest in obtaining atomic resolution images of various materials. However, the preparation of specimens for atomic resolution transmission electron microscopy is demanding, requiring a final specimen which is very thin (i.e.,  less than 50 nm) and free from artifacts. Typically, specimen preparation involves initial slicing, sectioning, trepanning, and/or grinding operations to produce a relatively thin (100-200 xcexcm) disk of approximately 3 mm diameter.\nIon beam milling systems have been used to prepare specimens of various materials including ceramics, semiconductors, metals, and combinations thereof for atomic resolution transmission electron microscopy. In such ion beam milling systems, such as the system disclosed in commonly-assigned U.S. Pat. No. 5,009,743, to Swann, specimens are mounted on holders and placed in the path of one or more ion beams. The ion beams gradually remove atoms from the surface of the specimen until a small perforation is formed in the center of the specimen. Generally, the area immediately around the perforation is then thin enough (i.e.,  less than 50 nm) for atomic resolution analysis by a transmission electron microscope.\nSwann et al, U.S. Pat. No. 5,472,566, provides a specimen holder adapted to permit the simultaneous two-sided ion beam milling of a specimen at very low angles of beam incidence, down to 0xc2x0, from both sides of the specimen through the use of one or more support arms to secure the specimen. This produces both rapid milling as well as reducing artifacts to provide high quality specimens for transmission electron microscopy analysis.\nIn such ion beam milling systems, the ion milled specimens have to be removed from the specimen holders and loaded in a TEM holder for imaging. Clearly, loading and unloading of the specimen from the holders increases the potential of damage to fragile specimens. Additionally, if the specimen is not completely electron transparent upon initial TEM examination, it must be returned to the ion mill for further thinning. Multiple milling steps may be required, depending upon thickness requirements at the area of interest in the specimen.\nSlope cutting by directing an ion beam at a masking screen located at or near a specimen surface is also known. The masking screen protects a portion of the specimen, while the remainder of the specimen is milled by the ion beam to reveal desired profiles. This technique is capable of providing cross-sectional profiles of layered structures, the surfaces of which can then be studied by scanning electron microscopy. For example, Hauffe, East German published Patent Application No. 201,538, teaches an apparatus and method for preparing specimens for SEM analysis by providing an ion gun inside the specimen chamber of a scanning electron microscope. A mask is positioned adjacent the sample and is used, in conjunction with the ion gun, to remove material from the specimen.\nDouble slope cutting to produce a very thin film of a specimen for TEM analysis has also been developed. Yoshioka et al, U.S. Pat. No. 5,907,157, describe a method and apparatus for preparing a specimen adapted for electron microscopy which includes an evacuated specimen-processing chamber. A specimen having a surface to be processed is placed inside the processing chamber and a beam-blocking member is placed close to the processed surface so as to block a part of an etching beam. A first etching step is performed by directing the beam at the specimen via the blocking member. Then, the specimen and the blocking member are moved relative to each other to expose another portion of the specimen. A second etching step is then performed by directing the beam at the specimen via the blocking member. As a result, the specimen is thinned and can be observed with an electron microscope. However, the requirement that the specimen and blocking member be moved relative to one another to perform the second etching step complicates the process and introduces the potential for errors as the relative positions of the two must be controlled to within an accuracy on the order of a few hundred nanometers (10xe2x88x929 meters) or less.\nMitro et al, U.S. Pat. No. 5,922,179, teaches an apparatus and process for etching and coating specimens in a single vacuum chamber, to minimize handling and transfer. The apparatus includes a sealed chamber and a vacuum pump for forming and maintaining a vacuum in the chamber, a first ion gun positioned in the chamber to etch a specimen, a sputtering target in the chamber, and at least one additional ion gun positioned in the chamber to cause material from the target to be directed onto the specimen.\nTechniques have also been developed in the preparation of specimens for scanning electron microscopy in which the specimen may be observed by a SEM during preparation thereof. For example, Grxc3xcnewald, U.S. Pat. No. 5,986,264, teaches an ion beam milling system using two ion guns for use in the preparation of samples to be used in both SEM and TEM analysis. The system uses an SEM as a high resolution imaging device to observe the progress of the ion milling of a sample and to determine when a proper sample thickness has been achieved. Also, as the sample stage can be tilted, the system is also taught to be useful for slope cutting of specimens.\nThe inadequacy in localized thinning of TEM specimens with broad ion beams, coupled with the ability of focused ion beam (FIB) to micro-machine materials, has led to the development of FIB milling systems. The first FIB based instrument (Precision Ion Milling System (PIMS)), became commercially available in 1984 from Gatan, Inc. Since the introduction of this instrument, several other commercial FIB systems have become available. These machines are capable of both imaging and milling selected areas within TEM and SEM specimens with submicron resolutions. The PIMS system is a gas ion source (argon) instrument, whereas the newer generation FIB systems employ liquid metal (gallium) ion sources having superior milling and imaging capabilities.\nSpecimens are normally milled in gallium based FIB systems with higher energy ions (compared to the PIMS), which may produce a thick amorphous surface artifact (damaged layer). This damaged layer should be removed to produce high quality specimens. Post-FIB milling in broad ion beam systems for final milling and cleaning of such specimens is a common practice. Specifically, FIB milled specimens have been re-milled and cleaned in broad ion beam machines using the specimen holder described in the above-mentioned U.S. Pat. No. 5,472,566 to improve specimen quality. However, again, this requires moving the specimen among various specimen preparation systems and the final TEM or SEM.\nOthers have reported the use of a wire shadow technique for TEM specimen preparation. Senz et al, Ultramicroscopy 70 (1997) pp. 23-28, and Langer, Elektronenoptische Untersuchung des Werkstoffverhaltesn in mechanisch belasteten Mikrobauteilen (1996) teach gluing a wire onto a specimen surface and exposing the surface to ion milling by a Gatan Duo-Mill. The non-shadowed areas of the specimen are removed, and a thin area is produced in the wire shadow.\nHowever, a need still exists in this art for simplified methods and techniques for specimen preparation for both SEM and TEM analysis.\nThe present invention provides a method and system for the preparation of specimens for analysis by electron microscopy which utilizes ion beam milling, but without the need for multiple positioning, cutting, and/or handling steps. The system and process are especially useful for the preparation for analysis by either TEM or SEM of semiconductors, metals, alloys, ceramics, and other inorganic materials.\nIn accordance with one aspect of the present invention, a system and process are provided for the preparation of specimens for analysis by transmission electron microscopy. The system includes a specimen processing chamber, at least two ion beam generators, and a specimen support or holder. Preferably, the system also includes the ability to view the progress of the milling operation and may include an imaging device such as a light microscope which permits monitoring of the area of interest on a specimen as the specimen is milled.\nSpecimens are prepared by initially forming, cutting, and/or polishing the specimen into a piece having a preferred maximum length of about 3 mm, a width of from about 500-1000 microns, and a thickness of from about 20-150 microns. The edge (thickness dimension and surface of interest for analysis) of the specimen should be relatively flat. A preferred shape for the specimen is a generally semi-circular disc having a radius of about 3 mm and a thickness of from about 20-150 microns, although other shapes may be utilized.\nAn ion beam masking member is then secured to the edge (thickness dimension and surface of interest for analysis) of the specimen using a suitable adhesive such as an epoxy glue or using a suitable mechanical attachment mechanism. The ion beam masking member is preferably a fiber having a diameter of less than about 100 microns, preferably less than about 80 microns, and most preferably less than about 35 microns. The fiber preferably has a length which approximates the length of the edge of the specimen, although the length of the fiber need only be as long as the area of interest on the specimen. Fibers having diameters down to about 7 microns have been found to be suitable for use. Alternatively, the ion beam masking member may have a non-circular cross-section including square, rectangular, elliptical, or other geometric shape. Preferably, the ion beam masking member comprises an amorphous or single crystal material such as carbon (either amorphous, graphitic, or diamond), silicon carbide, or sapphire. In one embodiment of the invention, the ion beam masking member may be pre-milled prior to securing it to the specimen through ion beam exposure to reduce its diameter or cross-section. This permits a more precise mounting and placement of the masking member on the surface of the specimen relative to features of interest on or in the specimen.\nFor TEM analysis, once the specimen has been mounted onto a support or holder and installed in the processing chamber, the specimen and ion beam masking member are exposed to two or more ion beams to mill and thin exposed areas on the edge of the specimen. Generally, the surface of interest on the specimen will be positioned perpendicular or nearly perpendicular to the direction of the ion beams. Preferably, an ion beam having a higher intensity in the center of the beam is utilized. That is, for an ion beam having a generally circular cross-section, the intensity at or near the center of the circle is higher than at the edges.\nThe ion beams will mill both the ion beam masking member as well as the specimen surface causing a necking or thinning of the masking member as milling is continued. Likewise, the masked area of the specimen beneath the ion beam masking member will also be thinned as milling continues because of the thinning of the masking member. The protected area beneath the masking member may be monitored using a suitable imaging device such as an optical or scanning electron microscope which is arranged generally perpendicular to the axis of the ion beams and which provides real time information concerning the progress of the milling process. Alternatively, a light emitting diode and photodetector may be arranged to provide an indication of when the specimen has been sufficiently thinned to become optically transparent. Additionally, if the ion beam masking member is electrically conductive, its electrical resistance may be measured and monitored such that a change in resistance is measured when the milling process causes the member to be severed (i.e., a gap forms).\nOnce the specimen has been thinned sufficiently to be electron transparent (i.e., less than about 200 nanometers, and preferably less than about 100 nanometers), it is removed from the processing chamber and is ready for TEM analysis. During milling, the specimen and ion beam masking member are positioned in tandem so that there is no relative movement between them. This avoids the need to stop the milling operation and accurately reposition the ion beam masking member prior to the resumption of milling as in some prior art processes.\nIn an alternative embodiment, the present invention provides a system and process for the preparation of specimens for analysis by a scanning electron microscope. In one form, the system includes a specimen processing chamber, at least one ion beam generator, and a specimen support or holder. Preferably, a second ion beam generator is also included. Generally, a specimen will be prepared so that it has a length of from about 4-8 mm, preferably about 6 mm; a height of from about 6 to about 12 mm, preferably about 8 mm; and a thickness of from about 250-750 microns, preferably about 500 microns.\nPreferably, the system also includes the ability to view the progress of the milling operation, in real time, and may include an imaging device such as a light microscope which permits monitoring of the area of interest on a specimen as the specimen is milled. Preferably, the system also includes a sputtering target which may be positioned so that one or both of the ion beam generators can sputter material from the target onto the specimen to coat its surface. Thus, once the specimen is etched and coated, it can be moved to an SEM for analysis.\nIn a preferred form, the specimen is secured to a rotatable specimen stage for milling. After milling has been completed, the specimen stage is rotated (e.g., 90xc2x0) to a position where the surface of the specimen is etched (cleaned). The specimen stage is then rotated back to its original position for sputter coating using one of the ion beam generators and a sputtering target.\nIn yet another embodiment of the present invention, the present invention provides a system and process for the preparation of specimens for analysis by a scanning electron microscope in a processing chamber which also includes an SEM. The SEM is arranged such that the specimen milling and preparation process is observable with the SEM being utilized as a high resolution imaging device. The ability to observe the milling process in real time permits precise positioning of the specimen and monitoring of the process. Further, the milled specimen can be analyzed directly in the chamber without movement to a separate apparatus.\nAccordingly, it is a feature of the present invention to provide a method and system for the preparation of specimens for analysis by electron microscopy which utilizes ion beam milling, but without the need for multiple cutting, positioning, and/or handling steps. It is another feature of the invention to provide a system and process which is adaptable for specimen preparation for both TEM and SEM analysis. These, and other features and advantages of the present invention, will become apparent from the following detailed description, the accompanying drawings, and the appended claims."}
{"text": "1. The Field of the Invention\nThe present invention relates to rendering graphical objects such as text. More specifically, the present invention relates to appropriately rendering a graphical object when a corresponding outline has excessive control points.\n2. Background and Related Art\nComputing technology has transformed the way we work and play. Computing systems now take a wide variety of forms including desktop computers, laptop computers, tablet PCs, Personal Digital Assistants (PDAs), and the like. Even household devices (such as refrigerators, ovens, sewing machines, security systems, and the like) have varying levels of processing capability and thus may be considered computing systems. As time moves forward, processing capability may be incorporated into a number of devices that traditionally did not have processing capability. Accordingly, the diversity of computing systems may likely increase.\nAlmost all computing systems that interface with human beings use a display to convey information. In many cases, the appeal of the display is considered an important attribute of the computing system. Display of textual information (e.g., Latin-based characters) typically includes processing glyphs that represent characters of a font. A glyph includes control points and instructions for connecting the control points such that an outline of a corresponding character can be generated in an arbitrary grid space (e.g., a pixel grid). Often, characters will be defined for display at a larger size and higher resolution and then mathematically scaled down (or otherwise manipulated) when the characters are to be rendered at smaller sizes and lower resolutions (or as bold, italic, etc.). Thus, a reduced number of descriptions, and potentially only one description, for a character (per font) need be stored.\nTo scale down a character, the location of control points can be divided by a scaling factor. For example, to scale a character down by a scaling factor of 10, the coordinates of each control point defining the character (at the higher resolution) can be divided by 10. It may be that control points defining a character for display on a 100×100 grid are to be scaled down for display on a 10×10 grid. Thus, a control point at grid position (50, 30) can be scaled down to a control point at grid position (5, 3), a control point at grid position (70, 70) can be scaled down to a control point at grid position (7, 7), etc. Accordingly, a smaller outline representing the character may be calculated and there is a reduced need for storing a number of different sizes of bit-maps for the character.\nThe smaller outline can then be analyzed to identify grid locations (e.g., pixels) that are to be turned on and that are to be turned off (a process often referred to as “scan conversion”). One scan conversion algorithm determines if the center of a grid position is inside or outside a resulting scaled down outline. When the center of a grid position is inside the scaled down outline the grid position is turned on. On the other hand, when the center of a grid position is outside the scaled down outline the grid position is turned off.\nUnfortunately, at times, and especially at lower resolutions, the results of scan conversion produce an unacceptable representation of a character. Unacceptable character representations can result from rounding errors in the scaling down process or from the scan conversion process itself. Further, rounding errors have a greater affect on character representation when a single grid location (e.g., a pixel) is on scale with the features of a character. For example, a rounding error that causes one grid location to be inappropriately turned on (or turned off) on a 50×50 grid may not even be detectable by the human eye. However, a rounding error that causes one grid location to be inappropriately turned on (or turned off) on a 4×4 grid may result in a character that is perceived as unacceptable to the human eye. Thus, at lower resolutions scan conversion can even fail to preserve the topology of the original outline of the character, for example, causing disconnected grid locations, inappropriately located features, and lack of symmetry.\nFor example, when an outline has a reduced number of control points associated with a corresponding character feature, scan conversion of the outline can result in the corresponding character feature being rendered at a non-optimal location in a grid space. A number of characters in Latin (Greek, Cyrillic, etc.) fonts have a horizontal stroke (“cross-bar”) somewhere between the top and the bottom of the character. For example, a capital “H” has a cross-bar approximately halfway between the top and the bottom. FIG. 1 depicts an outline 100 of Times New Roman capital “H”. By design, the position of the cross-bar is marginally above the mathematical halfway point of the vertical stroke of the “H”. That is, distance 101 is less than distance 102.\nHowever, by the time the “H” is rendered at a smaller size and lower resolution (e.g., 9 pt on a 96 dpi device), accumulated quantization errors can result in the cross-bar being rendered below the mathematical halfway point. Rendering the cross-bar below the half-way points is less than optimal, since the cross-bar was, by design, to be rendered above the halfway point. Pixel pattern 103 represents the pixels that would be turned on to render Times New Roman capital “H” at 9 pt and at 96 dpi resolution (i.e., the results of performing scan conversion on outline 100). As depicted in pixel pattern 103, the cross-bar is below the mathematical halfway point of the vertical stroke of the “H”.\nOne cause of non-optimal rendering results from the “H” having control points on the edges of the cross-bar, but not having control points, for example, on an imaginary line (represented by dashed line 104) along the center between the edges of the cross-bar. Traditional hinting solutions thus add a first constraint to a first control point, for example, on the top edge of the cross-bar (control point 107). The first constraint enforces a proportional position of the first control point in vertical direction relative to other control points on the base line (e.g. control point 109) and the caps line (e.g., control point 106). Subsequently, traditional solutions also add a second constraint to a second control point, for example, on the bottom edge of the cross-bar (control point 108). The second constraint enforces a distance between the second control point on the bottom edge and the previously constrained first control point (control point 106) on the top edge.\nSince both the first and second constraints will have to round positions in the vertical direction to the nearest grid position (e.g., pixel), each of the first and second constraints can introduce a rounding error of up to half a pixel, either way (up or down). In unfortunate cases, these rounding errors accumulate (e.g., both the first and second constraints round down)—to place the cross-bar below the mathematical halfway point.\nFurther, when an outline passes through the center of one or more grid locations, scan conversion of the outline can result in the character being rendered asymmetrically in a grid space. Italic fonts pose a particularly challenging problem because the jagged nature of diagonal strokes makes italic fonts more difficult to process, and quantization errors can compound the problem.\nAs previously mentioned, all scan-converters have to define is an interior pixel, even in cases where a pixel could just as well be an interior or an exterior pixel (e.g., when an outline passes through the center of the pixel). For example, TrueType® defines that when an outline is through the center of a pixel, the pixel is to be turned on. Thus, if an outline is a straight line, and if the line's run and rise (slope) reduce to a pair of odd numbers, the outline will pass through one or more pixel centers, and TrueType® will turn on a pixel, which may look like an extra pixel.\nReferring now to FIG. 2, FIG. 2 depicts an Arial Italic Capital “H”. The right edge of the left italic stroke (from control point 211 to control point 214) may rise 5 pixels for every pixel it runs across—a total of 2 pixels over and 10 pixels up, which reduces to a pair of odd numbers 1 and 5—and we can readily observe extra pixels (pixels 231, 232, and 233). Further, even when pixels having their centers exactly on an outline are turned off problems can still occur. For example, in FIG. 2, if pixels whose centers are exactly on the outline were to be turned off (instead of on), there would be two missing pixels (pixels 234 and 236), yielding a discontinuity in the pixel pattern. One approach is to define pixels whose centers are exactly on an outline to alternate between interior and exterior pixels as edges alternate between left and right edges. While this could solve the problem for some italic strokes, such an alternating approach would also introduce new problems on round strokes. For example, such an alternating approach could lead to asymmetrical pixel patterns in characters such as, an “O”.\nFurther, any movement of control points to comply with constraints can cause control points to be moved to inexact locations. Utilizing a control point that is at an inexact location for further calculations can result in non-optimal rendering of a corresponding character in a grid space. Often, hinting control points can include aligning some control points with other control points. For example, still referring to FIG. 2, control points 212 and 213 may have to be aligned with control points 211 and 214 because control points 211 and 214 may have been moved as a result of complying with their own constraints (e.g., constrained to caps and base lines respectively). While the act of aligning is a mathematically exact concept, its implementation in terms of a font hinting language (e.g., TrueType®) likely introduces numerical inaccuracies, because the internal numerical precision of font hint processors is typically limited. For example, the font hinting language TrueType® uses 6 binary places to store fractional control point location data and thus every numerical result can be accurate to only 26th or 1/64 of a pixel. Accordingly, the result of aligning, for example, control point 212 with its parent points 211 and 214 will be accurate to the nearest 1/64 of a pixel only—at best.\nMore specifically, since the line from control point 214 to control point 211 is 2 pixels over and 10 pixels up (as depicted in pixel pattern 203), it will run 0.2 pixels over for every pixel up. Accordingly, it should be 1.2 pixels over at control point 3, because control point 3 is 6 pixels up, or in terms of 1/64 of a pixel, it should be 1 12.8/64 pixel over. Unfortunately, 12.8/64 is beyond the precision limit of most, if not all, font hint processors. That is, 12.8/64 simply cannot be represented, and has to be rounded one way or the other (e.g., up or down to the nearest 1/64). Hint processors are typically configured with a default rounding behavior, for example, to always truncate, always round down, or always round up. In FIG. 2, and seemingly by chance, control points of the cross-bar are rounded down, even though rounding up would be a better approximation of the actual value (12.8 is closer to 13 than to 12). This error appears to be harmless, because it led neither to a spurious pixel, nor to a missing pixel.\nConversely, for control point 215, which again is over by 1 12.8/64 pixel, rounding down did lead to a spurious pixel (pixel 235), while rounding up would have prevented it. Moreover, the generality of typical font hint processors can produce results that are off by a fraction of a pixel, irrespective of the direction of rounding. Thus, even if rounding were to be correct, alignment problems can still result.\nAdditionally, character outlines that have too many control points can result in non-optimal rendering of a corresponding character in a grid space. For example, erroneous pixels can potentially arise when the exact shapes of outlines have increased complexity. For example, FIG. 3 depicts a Palatino Italic “H”. The Palatino Italic “H” includes significantly more control points than the Arial Italic “H” depicted in FIG. 2. The edges of the strokes of the Palatino Italic “H” are slightly curved, and thus include extra control points to model these curves. As depicted in FIG. 3, these control points are positioned in a seemingly random fashion on outline 300. Applying traditional scan conversion to outline 300 would result in a non-optimal pixel pattern, such as, for example, pixel pattern 303.\nUnfortunately, some hint processors refer to control points by consecutive number and removing a control point can cause other control points to be renumbered. Thus, renumbering control points can cause existing hinting instructions to refer to incorrect control point numbers. For example, removing what may be viewed as extraneous control points on outline 300 could cause existing hints to be inappropriately applied to any remaining control points (or not be applied at all). Inappropriate application of hints to remaining control points can significantly degrade a corresponding pixel pattern.\nAccordingly, to compensate for the sometimes non-optimal or inappropriate results of scan conversion (e.g., at lower point sizes and lower resolutions), a human typographer programmatically controls the scan conversion process on an ad-hoc basis by manually adding delta-exception hints (e.g., TrueType® instructions) to character outlines. Delta-exception hints are hints which are specific to individual characters and specific to a particular type size. That is, for each type size (of a font) at which there is an unfortunate quantization problem, the human typographer has to add possibly a plurality of delta-exception instructions, as appropriate. Since delta-exception hints are so specific, delta-exception hints are typically stored along with the character the delta-exception hints will be applied to (e.g., in a font file). Moreover, the use of delta-exception hints is an iterative process requiring repeated visual inspection of a plurality of type sizes before and after the application of corresponding exception instructions, which is laborious and prone to introduction of inconsistencies.\nDue in part to the wide variety of different artistic and technical features in different fonts, delta-exception hints must also be tailored to an individual font. That is, delta-exception hints for one font typically can not be reused for other fonts (even for the same character). Thus, for each new font, a typographer is required to iteratively determine appropriate delta-exception hints. Accordingly, the resources needed to compensate for non-optimal or inappropriate results of scan conversion across a variety of different fonts can be quite large. Therefore, what would be advantageous are mechanisms for automating the optimization of pixel patterns."}
{"text": "In performing operations within a cased well, such as logging formation properties, it is important to know the exact location of the tool lowered into the well to perform the specific function. Measuring the depth of well logging tools is traditionally made from the surface, by measuring how much of the cable which supports the tools has been deployed. Both wireline tools and while-drilling tools rely on the same basic concept.\nThe depth of the tool string is commonly determined by passing the cable over a calibrated measurement wheel at the surface of the well. As the tool is deployed, the length of cable unspoiled into the well is monitored as an estimate of tool depth. Depth compensation for cable stretch may be attempted by calculating a theoretical stretch ration based upon cable length, elasticity and tool weight. Even with very elaborate compensation algorithms, however, the actual amount of cable stretching may vary over time and because of unforeseen and unmeasured interactions between the cable and tool string and the well bore (such as tool hang-ups and cable friction) and anomalies such as cable “bounce”. Deviated wells, in which the tool is pulled along the interior surface of the well casing, can present particular problems with variable and inconsistent cable loading, as the stool “sticks” and jumps along the well bore. Such problems are also encountered albeit to a lesser degree, in tubing-conveyed operations in which tubing length is measured by a wheel arranged to roll along the tubing as it is unspoiled. Even very small deployment length measurement error percentages and other discrepancies can result, with either type of deployment, in absolute tool positioning errors of several feet or more in a well of over a mile in depth, for example.\nAnother approach has been developed that measures the rotations of a set of calibrated wheels contacting the cable under a set force which generates enough friction to transmit any linear velocity of the cable to the perimeter of the wheels, allowing a direct measurement of the corresponding depth increment.\nYet another approach utilizes a pair of sensors located on or within the drill string along with a known reference point within the borehole. Given the distance between the two sensors is determined by the a relatively short piece of drill string, its stretch/compression effects are negligible compared to the stretch and compression observed on the total-length of drill sting, and can be used as a downhole depth gauge by observing that when the second sensor will reach a correlatable event already seen by the first sensor, the depth increment is the distance between the sensors.\nUse of downhole sensors is challenging and present even more difficulties in the case of completed wells, in which the casing effects the ability to run certain downhole sensors. For example, steel casing has historically been thought of as a barrier to electromagnetic measurements of the properties of the formation. The problems presented by conductive liners are described by Augustin et al., in “A Theoretical Study of Surface-to-Borehole Electromagnetic Logging in Cased Holes,” Geophysics, Vol. 54, No. 1 (1989); Uchida et al., in “Effect of a Steel Casing on Crosshole EM Measurements,” SEG Annual Meeting, Texas (1991); and Wu et al., in “Influence of Steel Casing on Electromagnetic Signals,” Geophysics, Vol. 59, No. 3 (1994). These prior art references show that coupling between a transmitter and a conductive liner is independent of the surrounding geological formation conductivity for a wide range of practical formation resistivities encountered in the field and that the magnetic field produced inside the conductive liner at a distance of a few meters or less from the transmitter depends only on the conductive liner properties and not on the formation properties.\nFIG. 1 shows typical equipment used in the measurement of geological formation 10 resistivity between two drill holes 12a and 12b using electromagnetic induction. A transmitter T is located in one borehole, while a receiver R is placed in another borehole. The transmitter T typically consists of a coil (not shown) having a multi-turn loop (which consists of NT turns of wire) wrapped around a magnetically permeable core (mu-metal, ferrite or other ferro-magnetic material) with a cross section, AT. The transmitter T may further comprise a capacitor (not shown) for tuning the frequency of the coil. When an alternating current, IT, at a frequency of f0 Hz passes through this multi-turn loop, a time varying magnetic moment, MT, is produced in the transmitter. This magnetic moment is defined as follows:MT=NTITAT  (1)\nThe magnetic moment MT can be detected by the receiver R as a magnetic field, B0. The transmitter T, receiver R, or both are typically disposed in boreholes (e.g., 12a and 12b) in the earth formation 10. In this case, the detected magnetic field, B0, is proportional to the magnetic moment of the transmitter, MT, and to a geological factor, k1, as follows:B0=k1MT  (2)\nThe geological factor, k1, is a function of the spatial location and orientation of a field component of the magnetic field, B0, with respect to the magnetic moment of the transmitter, MT.\nThe receiver R typically includes one or more antennas (not shown). Each antenna includes a multi-turn loop of wire wound around a core of magnetically permeable metal or ferrite. The changing magnetic field sensed by the receiver R creates an induced voltage in the receiver coil (not shown). This induced voltage (VR) is a function of the detected magnetic field (BR), the frequency (f0), the number of turns (NR) of wire in the receiver coil, the effective cross-sectional area of the coil (AR), and the effective permeability (Ï□R) of the coil. Thus, VR can be defined as follows:VR=Ï□f0BRNRARÏ□R  (3)\nWhile f0 and NR are known, the product, ARÏ□R, is difficult to calculate. In practice, these constants may be grouped together as kR and equation (3) may be simplified as:VR=kRBR  (4)wherekR=Ï□f0NRARÏ□R.Thus, instead of determining the product ARÏ□R, it is more convenient to determine kR according to the following procedures. First, the receiver coil is calibrated in a known field, at a known frequency. Then, the exact value for kR is derived from the magnetic field (BR) and the measured voltage (VR) according to the following equation:kR=BR/VR  (5)\nWhen this system is placed in a conducting geological formation, the time-varying magnetic field, B0, which is produced by the transmitter magnetic moment, produces a voltage in the geological formation, which in turn drives a current therein, L1. The current, L1, is proportional-to the conductivity of the geological formation and is generally concentric about the longitudinal axis of the borehole. The magnetic field proximate to the borehole results from a free space field, called the primary magnetic field, while the field resulting from current L1 is called the secondary magnetic field.\nThe current, L1, is typically out of phase with respect to the transmitter current, IT. At very low frequencies, where the inductive reactance is small, the current, L1, is proportional to dB/dt and is 90° out of phase with respect to IT. As the frequency increases, the inductive reactance increases and the phase of the induced current, L1, increases to be greater than 90°. The secondary magnetic field induced by current L1 also has a phase shift relative to the induced current L1 and so the total magnetic field as detected by receiver R is complex.\nThe complex magnetic field detected by receiver R may be separated into two components: a real component, IR, which is in-phase with the transmitter current, IT, and an imaginary (or quadrature) component, II, which is phase-shifted by 90°. The values of the real component, IR, and the quadrature component, II, of the magnetic field at a given frequency and geometrical configuration uniquely specify the electrical resistivity of a homogeneous formation pierced by the drill holes. In an inhomogeneous geological formation, however, the complex field is measured at a succession of points along the longitudinal axis of the receiver borehole for each of a succession of transmitter locations. The multiplicity of measurements thus obtained can then be used to determine the inhomogeneous resistivity between the holes.\nIn both cases, i.e., measuring homogeneous geological formation resistivity or measuring inhomogeneous geological formation resistivity, the measurements are typically made before extraction of hydrocarbons takes place. This is because the boreholes typically are cased with conductive liners (e.g., metallic casing; see 16a and 16b in FIG. 3) in order to preserve the physical integrity of the borehole during hydrocarbon extraction. The conductive tubular liners interfere with resistivity measurements and are difficult and costly to remove from the borehole once they are installed. As a result, prior art systems such as that shown in FIG. 1 are not suitable for analyzing hydrocarbon reservoirs once extraction of the hydrocarbons begins.\nThe net or effective moment, Meff, of a transmitter inside a conductive liner is dictated by the inductive coupling between the transmitter and the conductive liner. Physically, the resistivity of the conductive liner is very low and the inductance relatively high. This property results in a current of almost the same magnitude as that of the transmitter current being induced in the conductive liner. Lenz's Law predicts that the magnetic field generated by this induced current in the conductive liner will oppose the time-varying magnetic field produced by the transmitter current. Thus, the magnetic field generated by the transmitter is mostly cancelled out by the magnetic field generated by the conductive liner. As a result, the magnetic field external to the conductive liner is greatly reduced, and its magnitude is proportional to the difference in currents in the transmitter and the conductive liner. In effect, the conductive liner “shields” the transmitter from any receiver positioned outside of the conductive liner.\nAn analogous situation is present with respect to a receiver if it is surrounded by a conductive liner. The field to be detected induces currents concentric with the receiver coil whose sense is such as to reduce the field within the liner. The field to be detected is consequently highly attenuated and the measurement is highly influenced by the variations in attenuation caused by the variation in liner properties, and example of which is graphically demonstrated by the slope of curve 10 shown in FIG. 2. The situation is exacerbated if both the transmitter and the receiver are surrounded by conductive liners. Often the design criteria for a crosshole survey of a cased borehole reduces the signal to a level that is undetectable by standard receivers. Moreover, the variance in conductivity, permeability, and thickness along a longitudinal axis of a liner makes difficult determining the attenuation factor at any given point.\nThe attenuation due to a steel casing surrounding the transmitter for a homogeneous formation is essentially constant a few meters form the source. Since the attenuation is constant the ratio of the fields as described, and incorporated herein, in U.S. Pat. No. 6,294,917 B1, to Nichols, removes the effects of the casing. This constancy of the fields can also be removed by calculating a shift operator as described, and incorporated herein, in U.S. Pat. No. 6,393,363 B1 to Wilt and Nichols. However, for multiple cased boreholes, both the ratio and shift operator methods preferably utilize an extra monitor for the transmitted field."}
{"text": "In contrast to semiconductor production of silicon wafers, which are processed individually, the manufacture of devices such as micromechanical sensors is accomplished using relatively small substrates. In this context, the base substrate is present in the form of a small cylindrical steel element. In order for the manufacturing process and for application of a circuit produced in thin-film technology onto such substrates to be cost-effective, the substrates are processed in multiple fashion. The individual base substrates are therefore retained in large numbers on a workpiece carrier, and processed together.\nHigh-pressure sensors that are utilized in many automotive engineering systems and in automation technology are fabricated on such substrates. Applications include, for example, direct fuel injection, common-rail technology for diesel vehicles, electrohydraulic brake systems, vehicle dynamics control systems, and many more.\nPressure is sensed via the deflection of the sensor membrane, which is coated with a Wheatstone measurement bridge using thin-film technology. The circuit is patterned using a photolithographic manufacturing process. Contacting and passivation can also be performed using a production step of this kind. The quality of high-pressure sensors depends substantially on the accuracy with which the resistance features of the Wheatstone measurement bridge are positioned on the pressure membrane. The resistances must be centered as much as possible on the pressure membrane. This positional accuracy substantially influences the electrical properties of the measurement bridge, for example, the range of the signal furnished by the sensor.\nIt is believed that the photolithographic process of patterning the functional layer to form the resistances is therefore essential in the production of a high-pressure sensor. A preconditioning of the surface of the functional layer, application of a photosensitive resist onto the functional layer, an oven step to condition the resist for exposure, exposure of the resist through an exposure mask, development of the exposed features, an oven step to condition the developed resist structure for etching, mapping of the resist image into the functional layer using an etching process, and subsequent stripping of the photoresist mask from the completed etched functional structure, are usually performed in this context.\nThe location of the resistances may be determined in the exposure step. The substrates, previously diced and retained on a workpiece carrier, may then be exposed with sufficient positional accuracy relative to an exposure mask. In contrast to the exposure of silicon wafers in semiconductor production, in which the wafers can be registered as one piece with the exposure mask (e.g. a quartz mask) to be imaged, in the exposure of the previously diced plurality of substrates there may be a positional inaccuracy for each substrate with respect to the imaging exposure mask.\nIt is believed that the substrates are often handled in immovably bolted-down workpiece carriers that must each be individually registered with respect to the exposure mask. This registration operation is performed manually.\nFor example, in the projection exposure method of the Nagano Keiki Co. Ltd., a special workpiece carrier that is automatically positionable is provided in each case. This workpiece carrier is not identical for the subsequent processes."}
{"text": "Hitherto, various types of structures of a multilayer printed wiring board including a built-up portion in which a plurality of wiring patterns that are stacked with insulating layers therebetween are electrically connected through via holes provided in the insulating layers have been proposed. For example, in such a multilayer printed wiring board, when a semiconductor device mounted is switched between the on and off states at a high speed, switching noise may be generated, resulting in an instantaneous drop in the electric potential of a power supply line. In order to suppress such an instantaneous drop in the electric potential, connection of a capacitor portion between the power supply line and a ground line to perform decoupling has been proposed. In Patent Document 1, provision of a layered capacitor portion functioning as such a capacitor portion in a multilayer printed wiring board has been proposed.\nPatent Document 1: Japanese Unexamined Patent Application Publication No. 2001-68858"}
{"text": "The present invention relates generally to protective garments and, more particularly, to an improved, more comfortable, more flexible, protective garment.\nPrior to the present invention, technological developments of yarns for protective garments have centered around the use of Kevlar, which is the DuPont trademark for an aramid fiber, as the fiber used in yarns, which yarns are ultimately used in protective garments. By way of example, and not by way of limitation, aramid fibers have been used in yarns with the yarns thereafter knitted to make protective garments including bullet-resistant vests and protective gloves as exemplified by Byrnes, U.S. Pat. No. 3,883,898.\nIn addition to the use of an aramid yarn as aforesaid, aramid fibers have been utilized in combination with other materials to form a yarn which yarn may be thereafter knitted to form a protective garment such as a protective glove with increased slash or cut resistance. Examples of this concept may be found in Byrnes U.S. Pat. No. 4,004,295, and Byrnes et al. U.S. Pat. No. 4,384,449, each of which describes the use of the Kevlar aramid fiber in combination with wire; the first of these two patents discloses the use of an aramid fiber yarn together with a metal wire and the second of these two patents describes a composite yarn itself; the yarn including a core of aramid fiber plus flexible wire and a covering of aramid fiber.\nThere are, of course, certain recognized problems with the use of the aramid fiber as the constituent in a yarn to thereafter be utilized in protective garments. For example, aramid fibers weaken in water. Second, the aramid fiber has only a limited resistance to true abrasion. Third, ultraviolet light adversely affects the appearance of the aramid fiber and can cause discolorations in the aramid fiber, discolorations in the yarn and discoloration in the finished product.\nRecently, a new high-strength fiber has been announced by Allied Corporation. The fiber is an extended chain polyethylene, which is a polyolefin, and has been sold under the trademark of Spectra with two different fibers being marketed, Spectra 1000 and Spectra 900. We understand that the Spectra 1000 is a 1200 denier fiber and that Spectra 900 is 650 denier fiber."}
{"text": "(1) Field of the Invention\nThe present invention relates to a process for producing aromatic imide polymer hollow filaments. More particularly, the present invention relates to a process for producing aromatic imide polymer hollow filaments useful for separating two or more different gases, for example, hydrogen gas and carbon monoxide gas, from each other with an excellent efficiency or for concentrating the gases from a dope solution of a specific aromatic imide polymer dissolved in a specific phenolic solvent by a wet membrane-forming method.\n(2) Description of the Related Art\nVarious methods for producing aromatic imide polymer semipermeable membranes or hollow filaments from a dope solution of an aromatic imide polymer in a phenolic solvent are disclosed by, for example, U.S. Pat. No. 4,378,400 for H. Makino et al and U.S. Pat. No. 4,378,324 for H. Makino et al.\nHowever, the products of the methods disclosed by the above-mentioned publications are still not satisfactory in the gas-separating properties thereof. That is, the aromatic imide polymer membranes disclosed in the above-mentioned publications are used as a separating membrance for a hydrogen-carbon monoxide mixture gas. The membrance exhibits an unsatisfactory selective gas permeability (P.sub.H.sbsb.2 /P.sub.CO) of about 20 to 65 and an unsatisfactory hydrogen gas permeating rate (P.sub.H.sbsb.2) of about 1.5.times.10.sup.-5 cm.sup.3 /cm.sup.2.sec.cmHg or less.\nThe term \"gas permeating rate\" used herein is defined as follows. ##EQU1## wherein X respresents an amount (volume) in cm.sup.3 of the gas passed through a membrane A represents a surface area in cm.sup.2 of the membrane through which the gas passed. T represents a transmission time in sec. of the gas through the membrane and D represents a difference in pressue in cmHg between the gas-supply side and the opposite side of the membrane\nAccordingly, it is strongly desirable to provide a process for producing aromatic imide polymer hollow filaments which are useful as gas-separating hollow filaments having an excellent gas selective permeability and a superior hydrogen gas permeating rate (P.sub.H.sbsb.2)."}
{"text": "Relay technology has been a hot area of research in wireless communication community in recent years. It is widely accepted that relay or multi-hop networks will become essential for beyond 3G mobile radio systems due to the range problem that appears there. As higher carrier frequencies or center frequencies can be envisaged for future mobile radio communication systems, where the expected center frequencies range up to 5-10 GHz and with bandwidth requirements of up to 100 MHz that can be foreseen, a significantly increased pathloss and noise power level has to be expected, which translates into a significantly reduced area a base station can cover.\nUsing relaying technology, signals can be transmitted much farther away, R. Pabst, B. Walke, D. Schultz, P. Herhold, H. Yanikomeroglu, S. Mukherjee, H. Viswanathan, M. Lott, W. Zirwas, M. Dohler, H. Aghvami, D. Falconer, and G. Fettweis, “Relay-based deployment concepts for wireless and mobile broadband radio,” IEEE Commun. Mag., vol. 42, no. 9, pp. 80-89, September 2004, J. Zhao, I. Hammerstroem, M. Kuhn, A. Wittneben, M. Herdin, and G. Bauch, “Coverage analysis for cellular systems with multiple antennas using decode-and-forward relays,” in Proc. 65th IEEE Veh. Tech. Conf., Dublin, Ireland, Apr. 22-25, 2007. Furthermore relay signals provide additional diversity and improve the received signal quality, J. N. Laneman, D. N. Tse, and G. W. Wornell, “Cooperative diversity in wireless networks: Efficient protocols and outage behavior,” IEEE Trans. Inform. Theory, vol. 50, no. 12, pp. 3062-3080, December 2004.\nIn order to circumvent the introduction of a denser grid of base stations (BS), the basic idea is to introduce relay stations (RS), which forward data packets to a mobile station (MS) that is out of reach of the base station. Such relay stations can be realized utilizing additional dedicated infrastructure relay stations having fixed power supplies, or they could be built into other mobile stations. Two main concepts of relaying were identified in the past: Amplify-and-Forward (AF) and Decode-and-Forward (DF). While AF has the advantage of being transparent to modulation and coding since a sampled version of the received signal is stored and retransmitted by the relay station without performing any decoding, DF allows for a separate adaptation to both links and avoids also the effect of noise enhancement since DF means that the relay station decodes and re-encodes the signal.\nCurrent relays cannot transmit and receive signals using the same time and frequency channel (half-duplex relays). Half-duplex generally means that communication is possible in both directions, but only one direction at a time (not simultaneously). Hence, traditional relaying schemes suffer from the spectral efficiency loss due to two channel uses.\nRecently, a spectral efficient relaying scheme called two-way relaying has been proposed in B. Rankov and A. Wittneben, “Spectral efficient protocols for half-duplex fading relay channels,” IEEE J. Select. Areas Commun., vol. 25, no. 2, pp. 379-389, February 2007 and in P. Larsson, N. Johansson, and K.-E. Sunell, “Coded bidirectional relaying,” in IEEE Veh. Tech. Conf., vol. 2, Melbourne, Australia, May 7-10, 2006, pp. 851-855. In such a scheme, the source and destination terminals transmit simultaneously in a first time interval. After receiving data signals from the source and the destination terminals in the first time interval, the relay station terminal retransmits the combined or superimposed source and destination data signals in a second time interval. Since both the source and the destination know its own data, respectively, both sides can cancel the so called self-interference, i.e. the own useful data transmitted towards the relay station in the first time interval, and decode the useful data from the other side. Two-way relaying achieves bi-directional transmission between the source and destination in two time intervals or time slots by using half-duplex terminals. Thus it avoids the spectral efficiency loss due to the use of half-duplex relays. Up to now, all the research is focusing on canceling the self-interference instead of utilizing it.\nThe authors of “Spectral efficient protocols for half-duplex fading relay channels” used a so called superposition coding scheme to re-encode data symbols to be transmitted in the second time interval, while the authors of “Coded bidirectional relaying” applied an XOR-operation on bit level to re-encode the data symbols to be transmitted in the second time interval at the relay station. The comparison of the two re-encoding schemes has been made in I. Hammerstroem, M. Kuhn, C. Esli, J. Zhao, A. Wittneben, and G. Bauch, “MIMO two-way relaying with transmit CSI at the relay,” in Proc. SPAWC, Helsinki, Finland, Jun. 17-20, 2007 and T. J. Oechtering, I. Bjelakovic, C. Schnurr, and H. Boche, “Broadcast capacity region of two-phase bidirectional relaying,” March 2007, submitted to IEEE Transactions on Information Theory.\nMIMO (Multiple-Input-Multiple-Output) technology has been a major breakthrough in the area of wireless communication in recent years. The use of multiple antennas can provide significant increase in capacity, I. E. Telatar, “Capacity of multi-antenna Gaussian channels,” Europ. Trans. Telecommun., vol. 10, no. 6, pp. 585-595, November 1999. It is almost certain that most communication systems will be equipped with multiple antennas in the years to come. Accurate channel knowledge may be used in order to take the advantage the MIMO system can offer. However, estimating the channel in a MIMO system is difficult. Traditional channel estimation schemes transmit dedicated, predetermined pilot or training sequences comprising known data symbols on different transmit antennas. Thus the spectrum efficiency is reduced in the traditional channel estimation schemes.\nMIMO channel estimation has been a research topic along with MIMO transmission technology. One of the most popular and widely used approaches to estimate a MIMO channel is to transmit orthogonal training sequences from different transmit antennas. The received signal is correlated with the training sequences at the receiver. Based on the received data and the knowledge of the predetermined training sequences, the channel matrices can be calculated. In M. Biguesh and A. B. Gershman, “Training-based MIMO channel estimation: A study of estimator tradeoffs and optimal training signals,” IEEE Trans. Signal Processing, vol. 54, no. 3, pp. 884-893, March 2006, the authors considered linear least squares (LS) and minimum mean-squared-error (MMSE) channel estimation approaches and investigated the optimal choice of training sequences. They showed that transmitting orthogonal training sequences from different antennas is optimal for the LS approach.\nInstead of multiplexing the training sequences with data symbols, H. Zhu, B. Farhang-Boroujeny, and C. Schlegel, “Pilot embedding for joint channel estimation and data detection in MIMO communication systems,” IEEE Commun. Lett., vol. 7, no. 1, pp. 30-32, January 2003, investigated the performance of embedding the training sequences with data to estimate the channel. Unlike conventional methods where predetermined pilot symbols are time-multiplexed with useful information data symbols, a pilot-embedding method was proposed in H. Zhu, B. Farhang-Boroujeny, and C. Schlegel, “Pilot embedding for joint channel estimation and data detection in MIMO communication systems,” IEEE Commun. Lett., vol. 7, no. 1, pp. 30-32, January 2003, where low power level predetermined pilots are transmitted concurrently with the useful data, and they are used to obtain an initial estimate of the channel such that a turbo decoding process can be started. Soft information obtained from the turbo decoder is subsequently used to improve channel estimates.\nA decision directed iterative channel estimation method was proposed in X. Deng, A. M. Haimovich, and J. Garcia-Frias, “Decision directed iterative channel estimation for MIMO systems,” in Proc. IEEE Int. Conf. on Communications, vol. 4, Anchorage, Ak., May 11-15, 2003, pp. 2326-2329.\nAll the proposed channel estimation schemes rely on an initial estimation of the MIMO channel based on the transmission of pure predetermined training sequences."}
{"text": "In a conventional WiMax/WiBro system, it is necessary for a BS to transmit a schedule in a header of each frame so that users may receive or transmit data at some location. In general, this schedule should be decoded by any user even the user with worse channel quality in a cell. Thus, the schedule is processed with lower level of modulation and encoding scheme and retransmitted for several times."}
{"text": "High bandwidth, short reach interconnections between chips or devices using conventional input/output (I/O) interfaces require significant power and chip area; specifically, solutions to maintain signal transfer quality between said chips or devices significantly contributes to the interface power consumption and smaller chip area. As a result, these conventional I/O interfaces are not desirable for low power and/or small chip area circuits."}
{"text": "1. Field of the Invention\nThe present invention relates to a novel settling accelerator for the epoxidation of a cyclic, at least monounsaturated alkene.\n2. Description of the Background\nNumerous processes for the epoxidation of alkenes are known and it is possible to use a wide range of different reaction or catalyst systems. The epoxidation of alkenes in a homogeneous, liquid phase using organic hydroperoxides in the presence of catalysts based on molybdenum, tungsten or vanadium is carried out industrially. However, epoxide production is accompanied by formation of equivalent or even greater amounts of the alcohol corresponding to the hydroperoxide, and the need to utilize or recirculate this alcohol greatly restricts industrial use of this process. For this reason, alternative processes for the more direct oxidation (epoxidation) of alkenes are being developed to an increasing extent.\nOne such process is epoxidation by means of molecular oxygen using silver catalysts. However, this process has only been able to be employed successfully for ethene and has not been able to be applied analogously to the epoxidation of other alkenes of the interest (for example propene).\nAnother process for direct oxidation of alkenes to epoxides is reaction with hydrogen peroxide. The process has been proposed for various epoxidation reactions because of, in particular, the positive aspect of the use of the oxidant in significantly reducing environmental pollution. Since the activity of hydrogen peroxide toward alkenes is only low, sometimes even completely absent, it is necessary to use activating agents, usually organic acids such as formic acid, acetic acid, etc., in organic solvents to facilitate the oxidation reaction. The acids form peracids in the reaction medium, which represent the actual reactive epoxidation agent, in situ. These processes, too, do not appear to be particularly successful, because it is difficult to obtain the peracids and because of the instability of the epoxides in acidic media, as a result of which varying convenient process conditions are necessary.\nStill another method is the oxidation of alkenes by reaction with highly concentrated hydrogen peroxide in a homogeneous, i.e. exclusively organic, liquid phase in the presence of soluble catalyst systems based on elements of Groups 4, 5 and 6 of the Periodic Table (Ti, V, Mo, W) in combination with elements selected from the group consisting of Pb, Sn, As, Sb, Bi, Hg, and the like. Here too, the results of the process do not permit implementation as an industrial process. This is firstly because the reaction proceeds slowly, and secondly the preparation of the catalyst systems, which generally consist of very complicated organic metal compounds and additionally have to be soluble in the organic reaction medium, is complicated and expensive. Furthermore, the use of highly concentrated hydrogen peroxide ( greater than 70%) involves considerable safety risks which cannot easily be overcome in an economical manner.\nThese processes of the prior art clearly show that the oxidation of alkenes by means of hydrogen peroxide is self-contradictory because the best working conditions in respect of the catalyst system and hydrogen peroxide involve an aqueous, acidic medium while the factors of the oxidation reaction itself and the stability of the epoxide are favored in a neutral organic medium. For this reason, further processes for the epoxidation of alkenes using hydrogen peroxide have been developed, in which either an improved catalyst system based on TiO2/SiO2 in an aqueous phase with addition of primary or secondary alcohols (see EP 0 987 259 A1) or a two phase system containing a catalyst comprising tungstic acid, a quaternary ammonium salt and a phosphorous compound (for example DE 30 27 349) is used.\nIn the case of alkenes whose epoxides are not hydrolysis-labile and in which the olefinic double bond is not sterically hindered (e.g. cyclic at least monounsaturated alkenes), the known epoxidation using hydrogen peroxide and a tungsten catalyst is the most economical alternative.\nFor the epoxidation reaction by means of hydrogen peroxide to proceed sufficiently quickly, a phase transfer catalyst (for example, Aliquat(copyright) 336) is usually used in the case of very lipophilic alkenes (for example cyclododecene)(Angew. Chem. (1991), 103(12), 1706-9). However, the desired strong acceleration of the epoxidation by means of the phase transfer catalyst leads to the phases being significantly more difficult to separate after the reaction, because of emulsion formation; the corresponding settling times increase greatly. In addition, the organic phase usually remains very turbid after the separation. To achieve virtually complete phase separation, it is necessary to use either phase separators having a very large volume or suitable centrifuges.\nThe increased settling times in this process greatly reduce its attractiveness for continuous, industrial-scale use. In particular, the process can usually not be implemented at all in existing plants because of space problems caused by the need for larger phase separators. The use of centrifuges is of little interest in view of the power costs and the maintenance requirement required by moving parts.\nIt can be said quite generally that settling times of less than 2 minutes are industrially desirable. If, on the other hand, the settling times are more than 4 minutes, an industrial-scale continuous process is difficult to operate economically.\nDE 30 27 349 describes a process for the epoxidation of alkenes using hydrogen peroxide, a tungsten compound, a phosphorus compound and a phase transfer catalyst. In this process, solvents such as alkanes or cycloalkanes are absolutely necessary. These solvents are always added to the reaction mixture in large amounts and generally serve either to dissolve a solid material and thus allow it to react, or to improve the reaction conditions, for example to achieve better heat removal.\nHowever, the dilution of starting materials with non-reactive substances such as solvents is undesirable since, firstly, the dilution leads to a reduction in the space-time yield and, secondly, a further separation operation is necessary after the reaction.\nAccordingly, one object of the present invention is to provide a settling accelerator, i.e. a compound, which leads to industrially acceptable settling times of the heterogeneous catalyst-containing reaction mixture in the epoxidation of alkenes, which ensures a sufficiently high space-time yield of epoxide, so that this process using the settling accelerator makes it possible for such epoxides to be prepared on an industrial scale.\nBriefly, this object and other objects of the present invention as hereinafter will become more readily apparent can be attained by a method of improving settling times of catalyst in the epoxidation of a cyclic, at least monounsaturated alkene, comprising:\nepoxidizing a cyclic, at least monounsaturated alkene having from 8 to 20 carbon atoms in the ring in a reaction medium containing an oxidant and a catalyst system comprising at least one metal of Groups 4, 5 and 6 of the Periodic Table of the Elements, phosphoric acid and a phase transfer catalyst and a cyclic alkane having from 8 to 20 carbon atoms in the ring, which corresponds to the alkene reactant, as settling accelerator in the epoxidation reaction."}
{"text": "Hexagonal containers are an efficient solution to maximize the usable volume within a container, and ultimately the amount of boxed merchandise that can be delivered per shipment. Additionally, hexagonal containers offer the advantage of minimizing the quantity of packaging material (e.g., cardboard) needed to package a given amount of merchandise. While the packaging industry makes use of these advantages, the hexagonal boxes currently available may be improved upon to further minimize the quantity of packaging material and glue necessary to package a given amount of merchandise."}
{"text": "1. Field of the Invention\nThe invention relates to a mount for reducing vibrations transmitted from a tool to a carrier vehicle, such as a tractor. More specifically, the mount isolates a tractor from the vibrations of a vibratory plow, or similar tool.\n2. Description of the Prior Art\nIn the past, different approaches have been tried to minimize the transfer of vibrations from a tool to a carrier vehicle, such as a tractor. The vibrations generally tire out the tractor operator and cause the bolted parts of the tractor to come apart, necessitating re-tightening of these bolts. The wear on the tractor parts also increases. The vibrating tool, for example a plow, is generally mounted on a sub-frame having a vibrating means attached to the sub-frame. The sub-frame is held by vibration dampening means in a larger frame. Examples of this construction are shown in U.S. Pat. No. 5,526,590 (Palm et al.), U.S. Pat. No. 5,482,121 (Draney et al.), U.S. Pat. No. 3,746,100 (Hall et al.), U.S. Pat. No. 4,140,425 (Flippin) and U.S. Pat. No. 4,164,982 (Draney). These solutions all share the apparent disadvantage of excessive wear and tear on the vibration dampening means, due to the fact that the sub-frame is vibrating in all directions. This subjects the vibration dampening means to compression (tensile) and shear stress. Most vibration means are rubber pads (or similar material), and inherently much less suited to shear stress, as compared to tensile stress. The vibration dampening means therefor wear out excessively fast in these applications. Also, the efficiency of vibration damping is not high enough to properly dampen vibrations from reaching the tractor.\nA further attempt to provide a vibratory plow construction with vibration dampening means is disclosed in U.S. Pat. No. 4,867,607 (Johnson et al.). This document discloses a pivotable part of the vibratory plow carrying frame. The frame elements are separated by air cushions (or resilient dampening means of other types), arranged to dampen vertical movement of the frame parts. This construction necessitates the use of extra frame parts, to achieve the pivoting of the frame, thus adding weight, complexity and cost to the construction. If solid rubber vibration dampers are used, they are objected to shear stress, which is detrimental to the life of such dampers."}
{"text": "1. Field of the Invention\nThe present invention relates to a DC-DC voltage converter. More particularly, the present invention relates to a DC-DC voltage converter with an output voltage lower than a reference voltage.\n2. Description of Related Art\nAs electronic information technique advances, circuit designs of various electronic products become increasingly complex. For many circuit designs, cell voltage is converted to lower power supply voltage, for example mobile phone, PDA and other hand held equipment with power supplied by cell. For example, for portable computer, voltage converter is used in a common power supplier to serve as a basic voltage adjuster, so increasingly more DC-DC voltage converters are used in various electronic products.\nThe function of the DC-DC voltage converter is outputting an output voltage smaller than an input voltage, and stabilizing the voltage. Referring to FIG. 1, a well-known DC-DC voltage converter 100 is shown. During operation of the DC-DC voltage converter 100 as shown in FIG. 1, it receives a DC input voltage Vin provided by an alkaline cell or a lithium cell, and the DC input voltage Vin is converted to a lower fixed DC output voltage Vout. After the output voltage Vout is fed back by a feedback circuit 101, it is compared with a reference voltage Vref through a comparator 103, and a control signal is generated to a driving circuit 105 according to a comparing result. The driving circuit 105 sends a driving signal to a power circuit 107 according to the received control signal. The power circuit 107 is connected to the input voltage Vin and a voltage stabilizing circuit 109. The power circuit 107 controls conducted/non-conducted state between the input voltage Vin and the voltage stabilizing circuit 109 according to the received driving signal, and then outputs a stable output voltage Vout after being stabilized by the voltage stabilizing circuit 109.\nHowever, in most of the portable electronic products, alkaline cells or lithium cells are used as the electric source. The voltage of the cell may be increasingly reduced with operation time and using time, so in order to obtain a stable output voltage, under many situations, the DC-DC voltage converter adopts boost-buck two-stage voltage converter.\nIn Taiwan patent No. 1281305, a dual-input voltage converter is provided, which includes a comparator for comparing the output Vout and a reference voltage Vref to generate a signal to a controller, and the controller determines to switch a first voltage converting circuit or a second voltage converting circuit. Referring to FIG. 2, a schematic view of an embodiment of this patent is shown. The first voltage converting circuit is composed of transistors SW1 and SW2, an inductor L and a capacitor Co, and the second voltage converting circuit is composed of transistors SW2 and SW3, the inductor L, and the capacitor Co. The controller 201 switches the transistors SW1 and SW2 to generate a current I to charge the capacitor Co through the inductor L to obtain an output voltage Vout. The comparator 203 compares the output voltage Vout and the reference voltage Vref. When the output voltage Vout is lower than the reference voltage Vref, an output signal Se of the comparator 203 makes the controller 201 turn off the transistor SW1, such that a voltage VDD is converted to a voltage VPP after being boosted by a boost converter 205. The controller 201 switches the transistors SW2 and SW3 to convert the voltage VPP to the output voltage Vout, such that the dual-input voltage converter 200 can be considered as a boost-buck two-stage voltage converter, and avoids the situation that the output voltage Vout is reduced due to the reduce of the input voltage Vin and thus causing an unstable output. However, the current DC-DC converter still has some defects which are described as follows.\n1. Although the two-stage voltage converter can solve the situation that the reduction of the input voltage Vin results in the unstable output voltage Vout, it cannot operate under a state that the output voltage Vout is smaller than the reference voltage Vref, and it depends on the boost converter to make the output voltage Vout maintain at the state of being larger than the reference voltage Vref. However, low operating voltage enables the electronic product integrate more transistors on a single chip, so in practical application, various electronic products use low voltage increasingly. Through prediction of “International Technology Roadmap for Semiconductors (ITRS)” issued by Semiconductor Industry Association (SIA), in the period from now to the year of 2016, operating voltage and operating current of IC will respectively show falling and rising trends. It can be known that the output voltage Vout of the voltage converter required by many electronic products becomes increasingly smaller, so as to directly result in the difficulty of selecting value of the reference voltage Vref. A situation that the output voltage Vout is smaller than the reference voltage Vref occurs more frequently. It is impossible for the conventional DC-DC voltage converter to operate under the situation that the output voltage Vout is smaller than the reference voltage Vref, and it is a blind spot in the design of the DC-DC voltage converter.\n2. Under the situation that the output voltage Vout is smaller than the reference voltage Vref, the dual-input voltage converter adopts the method of boosting and then bucking. However, it is necessary to add a boost converter, a voltage converting circuit, and some other related electronic components in the circuit. However, PCB space of the current DC-DC voltage converter is quite limited, which does not satisfy the developing trend of the modern electronic product of “a small size with big function”. The DC-DC voltage converter is widely used in the electronic products, and the increase of the electronic components in the DC-DC voltage converter will result in the increase of manufacturing cost."}
{"text": "As is generally known, computers are used to manipulate data under the control of software. Software is typically written in a high level programming language such as C, which is then compiled by a compiler program into binary machine language instructions which can be executed by a central processing unit in the computer. Software written in programming languages other than the machine language instructions are relatively much easier to understand and use. A very common strategy to simplify computer programming is to group frequently called portions of a program in subroutines, or functions, which perform a certain task. Functions may generally be called or executed as needed from anywhere else in the program each time the task is to be performed. Thus, rather than repeating the program instructions for the task each time the task is performed, the program instructions appear only once in the software, reducing the size of the resulting software.\nHowever, including functions in software can have a negative impact on the performance of the software. Information must often be passed along to the functions as formal parameters, and functions may return information as return values. Passing formal parameters and return values to and from functions requires that a compiler generate additional machine language instructions, making extra work which slows the computer. Additionally, each time a function is called, the computer must save the state of the processor before jumping to the function in the program, in effect saving its place before executing the function. Once the function has finished executing, the computer must restore the state of the processor before returning to the instruction following the function call. These additional tasks can greatly slow the execution of software, particularly if the software includes many small functions. In the extreme, these additional tasks such as passing parameters and return values and saving the state of the processor can be more work for the processor than the actual function.\nMost compilers include optional optimization tools which give the programmer the option of optimizing the software for speed as it is compiled. One such tool, called inlining, replaces function calls with the body of the function each time the function call appears. Thus multiple copies of a function will appear inline with the rest of the program, rather than being a single independent copy which can be called from multiple places. Actual arguments to an operation are substituted for formal parameters and the state-saving instructions are omitted. Since this increases the size of the software, compilers typically attempt to determine which functions would be good candidates for inlining and which are not. The criteria used by compilers include information such as the size of a function, the frequency with which it is called, and the number of places from which it is called. For example, large functions which are called infrequently but from many points, or call sites, in a program would be poor candidates for inlining as they would greatly increase the size of the software. In contrast, small functions which are called frequently from only a few call sites would be good candidates for inlining.\nHowever, the traditional software environment in which a program is written, compiled, then executed on a single type of computer is changing with the increasing desire for hardware independent software which can be executed on multiple different types of computers. To achieve hardware independence in software, software is typically either manipulated after compilation (e.g., translated) or is compiled piecemeal as the software is executed.\nAnother departure from the traditional compilation environment is dynamic translation, wherein software in binary machine language form, written for execution on a first type of computer, is translated as it is executed on a second type of computer. Dynamic translators operate by translating each word of the non-native code into a corresponding word or words of native code for execution by the computer. However, dynamic translators do not scan, evaluate and modify software before executing, and thus do not have detailed information about the software prior to execution. Without additional information about the code it is very difficult, if not impossible, to achieve full inlining of the program to improve performance, particularly in machine language code.\nA need therefore exists for a method of enabling function inlining and related optimizations during execution of a program, whether native or non-native to the computer. A further need exists for a method of enabling function inlining and related optimizations during execution of a program in machine language binary format. A further need exists for a method of providing information about an executing program to facilitate dynamic optimization of the program."}
{"text": "The present invention relates to a process and apparatus for forming plastic sheet. In particular, the present invention relates to a process and apparatus for forming plastic sheet having low residual stress and high surface quality. Plastic sheet formed according to the process of the present invention is particularly useful in optical and electronic display applications, such as, for example, optical windows, optical filters, recording media, and liquid crystal displays (LCD).\nSheets of optical quality glass or quartz are used in electronic display applications as \"substrates.\" In such applications, a \"substrate\" is a sheet of material used to build an electronic display. Such substrates can be transparent, translucent or opaque, but are typically transparent. In general, such sheets have conductive coatings applied thereto prior to use as substrates. Such substrates often have stringent specifications for optical clarity, flatness and minimal birefringence, and typically must have high resistance to gas and solvent permeation. Mechanical properties such as flexibility, impact resistance, hardness and scratch resistance are also important considerations. Glass or quartz sheets have been used in display applications because these materials are able to meet the optical and flatness requirements and have good thermal and chemical resistance and barrier properties; however, these materials do not have some of the desired mechanical properties, most notably low density, flexibility and impact resistance.\nBecause of the mechanical limitations of glass or quartz sheet in optical or display applications, it is desirable to use plastic sheet in such applications. Although plastic sheets have greater flexibility, are more resistant to breakage, and are of lighter weight than glass or quartz sheets of equal thickness, it has been very difficult to produce plastic sheet having the requisite optical specifications needed for use in optical and display applications at reasonable costs. Moreover, many types of plastic sheet undergo unacceptable dimensional distortion when subjected to substrate processing conditions during manufacture of the display devices, particularly with respect to temperature.\nThere are several commercially utilized methods for producing plastic sheet, including casting, extrusion, molding, and stretching operations. Of these methods, several are not suitable for producing high quality plastic sheet. As used throughout this specification, the term \"high quality\" is used to describe plastic sheet having the following characteristics: low surface roughness, low waviness, low thickness variation, and minimal amount of polymer chain orientation (for example, as measured by asymmetric physical properties, birefringence or thermal shrinkage).\nFor example, injection molding is likely to produce high amounts of polymer chain orientation, especially for thin sheets (i.e., 1 mm thickness or less), due to the flow of molten plastic into the mold, which unacceptably increases birefringence for polymers with non-negligible photoelasticity coefficients. Injection compression molding is an improved molding process which allows squeezing of the polymer after injection for the purpose of improving surface quality and reducing polymer chain orientation. However, even with these improvements, injection compression molding has limited ability to produce high quality sheet.\nCompression molding and press polishing may be used to produce sheets with good surface quality; however, the squeezing flow inherent in such processes results in polymer chain orientation which results in unacceptable shrinkage during thermal cycling. Moreover, these processes are not continuously operable and therefore increase labor and production costs.\nStretching operations (for example, for the production of uniaxially- or biaxially-oriented films) and blown film extrusion inherently introduce large amounts of polymer chain orientation and are unsuited for the production of high quality plastic sheet.\nSolvent casting can be used to produce high quality film; however, there are practical limitations to the maximum film thickness which can be produced by this method. In addition, the solvent used in the casting must be removed after formation of the sheet.\nSheet extrusion is run as a continuous operation, but this process introduces unacceptable polymer chain orientation due to the nature of the polymer flow in the die and between the polished rollers in the roll stack.\nThere is therefore a continuing need for a method for producing relatively inexpensive, high quality plastic sheet in a continuous fashion, wherein the resultant plastic sheet is capable of use as a substrate in optical and electronic display applications."}
{"text": "The present invention relates to a roller formed of a sintered alloy having high wear resistance and assembled in a rotary compressor having high fluid tightness.\nIn a recent trend, a rotary comprressor for use with domestic electrifications becomes light in weight and more compact in size. Further, for the reduction of production cost and for high performance of the compressor, improvements have been made on materials of respective mechanical components of the compressor.\nA rotary compressor mainly includes, as shown in FIGS. 1 and 2, an outer case 10, a cylindrical housing 11 assembled in the case 10 and formed with a vane groove 11A extending in radial direction of the housing, a roller 13 rotatable eccentrically in the housing 11, a shaft 14 integrally fixed to the roller 13 for its rotation, and a vane 12 slidably disposed in a vane groove 11A and moved in radial direction of the roller 13. A compression spring 15 is disposed in the groove 11A to urge the vane 12 radially inwardly. Therefore, a radially inner end face of the vane 12 is in slide contact with an outer peripheral surface of the roller 13. A working chamber is provided by a space defined between the housing 11 and the roller 13, and the vane 12 divides the chamber into intake and discharge chambers. The intake chamber is connected to an intake port 16 and the discharge chamber is connected to a discharge port 17. A fluid sucked in the working chamber is compressed and fed out by the eccentric rotation of the roller 13.\nAmong those components, the vane 12 and roller 13 perform relative sliding motion at high load, and therefore, these components must have high wear resistances. On this standpoint, various materials, for example, sintered alloy, have been proposed for the materials of the vane and roller.\nHowever, regarding the material of the vane, SKH51 has been still a major material in an actual production. SKH 51 is defined by JIS (Japanese Industrial Standard), which is a high-speed tool steel, containing 0.80 to 0.90% of C, 3.80 to 4.50% of Cr, 4.50 to 5.50% of Mo, 5.50 to 6.70% of W, 1.60 to 2.20% of V and balance Fe.\nFurther, regarding the material of the roller, a sintered material has been employed as described above rather than a cast iron. According to the proposed sintered material, a hard metal carbide and a metal oxide formed by a steam treatment are dispersed in a matrix. Such sintered alloy is disclosed in Japanese Patent Application Kokai Nos. 60-73082 and 60-174853.\nMore specifically, according to the publication No. 60-174853, it discloses a sintered alloy consisting of 3-10% by weight of chromium, 1-5% by weight of graphite and balance iron and impurities. Metal carbide, metal oxide and free graphite are dispersed in the tempered martensitic matrix. The metal oxide is formed at interiors of sintered voids by steam treatment. This metal oxide seals the sintered voids to thereby obtain lubrication oil retainability.\nAccording to the publication No. 60-73082, it discloses a rotary compressor. A rotor and/or a vane are formed of ferrous sintered alloy. Metal carbide and metal oxide are formed during tempering and are dispersed in the tempered martensitic matrix. Further, nitrogen is solid-solved in the martensitic matrix.\nThe above sintered material for the roller is intended to improve wear resistance and fluid-tightness. The metal oxide which seals sintered pores serves to enhance fluid-tightness of the compressor. However, with the use of such sintered material, a compressor roller has been burdened with much higher load because of the recent use of an inverter system. Accordingly, such sintered material may be worn out and undergoes scuffing if applied in the inverter system. In order to overcome this problem, there is a further demand for further improvement on wear resistivity and scuffing by using specific composition instead of conventional metal carbide dispersed in the matrix."}
{"text": "This invention relates to an accumulation system for driverless vehicles of the type which are propelled along a track by frictional contact between a drive wheel and a rotating drive shaft. In particular, this invention relates to the deceleration and accumulation of a number of driverless vehicles along a track without the need for direct cooperation between one vehicle and a fixed cam on the rear end of an adjacent vehicle. This invention may also be utilized to cause the vehicle to reverse its direction of motion along the track.\nAn accumulation system for driverless vehicles is shown in U.S. Pat. No. 3,818,837 and includes a speed control device or tail which is attached to the rear end of the vehicle for cooperation with a drive control member on the front end of the next adjacent vehicle. The speed control device is typically in the form of a cam. The contact between the tail on one vehicle and the drive control member on an adjacent vehicle on the track causes the trailing vehicle to accumulate or stop adjacent to the forward vehicle. Release of the forward vehicle will also cause a corresponding release of the trailing vehicle. The cam on the tail is contoured such that the change in speed of the trailing vehicle is effectuated gradually."}
{"text": "With the development of intelligent mobile terminals and the popularity of social communication networks, the social communication scope of an intelligent mobile terminal user is greatly expanded. Therefore, the mobile terminal user has lots of contact information to manage. The current contact management manner has some obvious disadvantages in managing large amounts of contact information. For example, a contact is difficult to search for or find because the contact is generally searched for by using the first letter of an alphabetic identifier of the contact, that is, inputting the first letter of the contact in a contact search box. Because the contact information is huge, a large number of contacts having the same first letter may be searched out or found by using the first letter. These contacts having the same first letter are all displayed. In this case, the user is only capable of finding needed contact information by browsing or scrolling pages. This process is time-consuming, and does not satisfy a user's expectation for quickly searching for and locating a contact. In addition, user experience on the intelligent mobile terminal is reduced."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a microwave oven of the type including a duty cycle control which periodically energizes and de-energizes the microwave generating system during the cooking process and, more particularly, to a circuit arrangement for providing an improved delay start feature in such an oven without costly changes in the programming or circuitry which controls oven operation.\n2. Description of the Prior Art\nIt is common in commercially available, electronically controlled microwave ovens to provide a delay start feature. Such a feature permits a microwave oven user to program the control circuitry of the oven so that a cooking operation begins at a later time, the delay time being entered by the user into the oven controls. Thus, an operator of an oven desiring to use such a feature places food in the oven; selects a particular power level at which he desires to cook the food and the time duration for the food to cook and enters this information into the oven control system; selects and enters a time delay (usually in the form of a preselected number of minutes) into the control, then usually presses a \"start\" button on the control signalling the completion of cooking instructions, and leaves the oven unattended.\nMicrowave ovens heat food by subjecting it to microwave radiation generated by a magnetron. Most microwave ovens vary the cooking power level by using a duty cycle control to operate the magnetron. In the operation of a duty cycle control, the magnetron is switched between a full-on condition and a full-off condition with the percentage of \"on\" time compared to the total time of each timing period (duty cycle) being varied to change the cooking power delivered to the food. Various specific circuits to provide this duty cycle are well known in the art. These range from simple cam operated mechanical timers having electrical controls to operate the magnetron to more sophisticated systems employing electronic solid state timing and switching elements. Since the duty cycle control does not form a part of this invention, it will not be described in detail. The magnetron therefore represents a first class of load in a microwave oven which is energized intermittently during the cooking process.\nAs contrasted to the intermittent operation of the magnetron, the microwave oven also includes various other electrical loads which are desirably energized continuously during the cooking process, such as the mode stirrer motor, the magnetron cooling blower and the filament heater for the magnetron.\nIn microwave ovens having a delay start feature, provision is usually made to control the continuously energized loads separately from the magnetron. More specifically, it is not possible to control these continuously energized loads using the same switching device that controls the magnetron because the latter is operated in an intermittent or duty cycle controlled fashion.\nOne solution to this problem is to turn on the constantly energized loads when the oven user completes entry of the delay start instructions into the oven controls. In this arrangement, the mode stirrer motor, blower, etc., are energized by activating the start button, or other equivalent input control device and remain energized continuously during the delay time period when the magnetron is idle, waiting for the beginning of the cooking cycle.\nThe obvious disadvantage of such a scheme is that power used to drive these loads is wasted during the time when they are not functionally needed to support the cooking operation. In addition, the noise caused by operation of these loads is undesirable and distracting.\nAnother solution to the above mentioned problem is to control the continuously energized loads by means of a separate independent signal developed by the oven control after the time delay period has passed at the beginning of the cooking cycle, i.e. when the magnetron is initially energized. This may be typically done by utilizing an additional input/output port on the microprocessor control with appropriate programming to carry out this function. In addition, an electronic switching element, associated drive electronics and isolators must be employed to implement this approach, resulting in a significantly additional cost. Moreover, it is sometimes difficult to couple additional electronics to an already existing control package without a significant redesign effort."}
{"text": "Organizations such as on-line retailers, Internet service providers, search providers, financial institutions, universities, and other computing-intensive organizations often conduct computer operations from large scale computing facilities. Such computing facilities house and accommodate a large amount of server, network, and computer equipment to process, store, and exchange data as needed to carry out an organization's operations. Typically, a computer room of a computing facility includes many server racks. Each server rack, in turn, includes many servers and associated computer equipment.\nSome computer systems typically include a number of components that generate waste heat. Such components include printed circuit boards, mass storage devices, power supplies, and processors. For example, some computers with multiple processors may generate 250 watts of waste heat. Some known computer systems include a plurality of such larger, multiple-processor computers that are configured into rack-mounted components, and then are subsequently positioned within a rack system. Some known rack systems include 40 such rack-mounted components and such rack systems will therefore generate as much as 10 kilowatts of waste heat. Moreover, some known data centers include a plurality of such rack systems.\nSome computer systems, which may function as servers, include a number of components that are mounted in an interior of the computer systems. The components, which can include printed circuit boards (for example, a motherboard) and mass storage devices, can support one or more components that generate waste heat, referred to herein as “heat-producing components.” For example, a motherboard can support a central processing unit, and mass storage devices can include hard disk drives which include motors and electronic components that generate heat. Mass storage devices can also include solid state drives that generate heat. Some or all of this heat must be removed from the components to maintain continuous operation of a computer system. The amount of heat generated by the central processing units, hard disk drives, etc. within a data room may be substantial. Heat may be removed from the heat-producing components via an airflow flowing through a computer system.\nIn some cases, cooling systems, including air moving systems, may be used to induce airflow through one or more portions of a data center, including airflow through a rack-mounted server that includes various heat-producing components. However, some servers direct airflow through an interior that includes multiple heat-producing components. Air removes heat from the heat-producing components and in turn the air is heated as it passes through the interiors that include the various heat-producing components, so that air passing over heat-producing components in a downstream portion of the server has a higher temperature than air passing over heat-producing components in an upstream portion of the server. Higher temperatures of the air in the downstream portion of the server may result in reduced heat removal capacity relative to air passing over heat-producing components in the upstream portion of the server. As a result, less heat may be removed from downstream heat-producing components than upstream heat-producing components. In some cases, the downstream heat-producing components are more sensitive to heat than the upstream heat-producing components, which can result in a suboptimal configuration.\nIn some cases, a computer system mounted in a rack includes one or more hot-pluggable electronic components, which can be added, removed, swapped, etc. from a computer system without powering down the computer system. Hot-pluggable electronic components in a computer system are often mounted at an external side of the computer system, including a “front” side through which cooling air is received into the computer system, to enable simplified access to the components for removal, addition, swapping, etc. In addition, mounting hot-pluggable electronic components to an external side of the computer system enables hot swapping without moving the computer system itself. Mounting hot-pluggable electronic components throughout the depth of the computer system interior can hamper hot-swapping efforts while maintaining operations by other hot-pluggable electronic components. Furthermore, as indicated above, mounting hot-pluggable electronic components throughout the depth of the interior can result in preheating of cooling air which removes heat from some components which are downstream of other components, which can reduce cooling efficiency and can negatively affect component performance.\nThe various embodiments described herein are susceptible to various modifications and alternative forms. Specific embodiments are shown by way of example in the drawings and will herein be described in detail. It should be understood, however, that the drawings and detailed description thereto are not intended to limit the disclosure to the particular form disclosed, but on the contrary, the intention is to cover all modifications, equivalents and alternatives falling within the spirit and scope of the appended claims. The headings used herein are for organizational purposes only and are not meant to be used to limit the scope of the description or the claims. As used throughout this application, the word “may” is used in a permissive sense (i.e., meaning having the potential to), rather than the mandatory sense (i.e., meaning must). Similarly, the words “include,” “including,” and “includes” mean including, but not limited to."}
{"text": "α2-Antiplasmin (α2AP) is a glycoprotein in blood plasma that rapidly and specifically inhibits the enzyme, plasmin, which digests blood clots, whether presenting early as intravascular platelet-fibrin deposits or as partially or completely occlusive thrombi. Similarly, plasmin and α2AP activities are important to the development and survival of fibrin as occurs in inflammation, wound healing and virtually all forms of cancer and its metastases. Human α2-antiplasmin (α2AP), also known as α2-plasmin inhibitor, is the main inhibitor of plasmin. Plasmin plays a critical role in fibrin proteolysis and tissue remodeling. The inventors discovered antiplasmin-cleaving enzyme (APCE) in human plasma and showed that it is a soluble isoform or derivative of fibroblast activation protein-alpha (FAP). Like APCE, FAP is also a prolyl-specific enzyme that exhibits both endopeptidase and dipeptidyl peptidase activities.\nFAP significantly over-expresses in >90% of epithelial-derived cancers [1-3]. FAP is produced transiently by activated stromal fibroblasts during embryogenesis [4] and wound healing [3], but other than an occasional normal fibroblast or pancreatic islet α-cell, it is not expressed by normal adult tissues or benign tumors [2, 3, 5]. FAP is prominent on the membranes of proliferating fibroblasts in diseases where fibrous tissue growth is a conspicuous feature, such as primary pulmonary fibrosis [6]; chronic hepatitis [7]; certain bone-associated malignancies [8, 9]; and the arthritides [10]. Selected parenchymatous cancer cells may also occasionally express FAP [11].\nWhile a biologic substrate for the proteinase activity of FAP has not been definitively established, results indicate that FAP helps digest extracellular matrix (ECM) components as tissue is remodeled to accommodate cancer expansion [2, 16, 17]. Paradoxically, activated fibroblasts not only digest ECM, but also synthesize ECM components of the stromal scaffolding that support cell division and motility during neoplastic growth [18]. FAP has been considered a potential target in the diagnosis and therapy of cancer, with inhibition of FAP proteinase activity posed as possibly therapeutically useful [19, 20]. Santos et al. [21] have shown that genetic deletion or pharmacologic inhibition of FAP by glutamyl-proline boronic acid (Glu-boroPro) decreased stromal growth in mouse models of lung and colon cancer. However, Glu-boroPro has an exceptionally short plasma half-life before cyclizing and losing inhibitory activity [22]. Moreover, it also inhibits dipeptidyl peptidase IV (DPPIV), which is important in plasma glucose regulation and immune function [23]. Hence, despite inhibiting FAP and suppressing tumor growth, Glu-boroPro is not likely to be therapeutically useful in cancer [24].\nThe measurement of cellular FAP activity and inhibition is confounded by another prolyl endopeptidase: namely, Prolyl Oligopeptidase (POP) which is elevated in many cancers [25]. POP is a prolyl-specific serine proteinases, which cleaves peptides of less than about 30 amino acids in length. The enzyme is present in most tissues, but is noted to be more abundant in selected organs, e.g., brain and kidney. Recently POP has been indicated as making secondary cleavages of thymosin-β4 to yield the derivative peptide, acetyl-serine-asparagine-lysine-proline, which appears to be a potent stimulator of angiogenesis [26]. Both FAP and POP activities are commonly measured using non-specific substrates such as Z-Gly-Pro-AMC or succinyl-Gly-Pro-AMC, neither of which distinguishes between the two activities [27]. Consequently, total prolyl-specific endopeptidase activity, which is often attributed to FAP alone, may also include POP activity. This complicates interpretations about the effects of inhibiting either enzyme on cancer growth. Pre-clinical studies have suggested other promising applications of POP inhibition for managing memory, learning disorders and depression, but development of relatively benign, highly effective POP inhibitors for in vivo testing have been elusive. The results, newly described herein, indicate that POP is expressed in significant amounts by a variety of cancer cells grown in culture.\nCertain compounds which specifically inhibit either one or both of FAP and POP and therefore can be used to treat various conditions which involve these proteins are desirable. Thus, the presently disclosed and claimed inventive concept(s) is directed to, but not limited to, substrates and inhibitors of FAP and/or POP, and to methods of using FAP and/or POP inhibitors for treating conditions such as but not limited to, cancers, neural disorders, and angiogenesis, and to screening methods for identifying such inhibitors, that overcome the disadvantages and defects of the prior art."}
{"text": "Consumption of nutrients, like antioxidants and folic acid, which are abundant in fruits and vegetables, has been linked to a lower incidence of cardiovascular disease. Moreover, it is well settled that eating fruits high in soluble fiber can reduce cholesterol levels which protects against atherosclerosis.\nOther advantages of having a diet high in fruit include better athletic performances, reduced risk of developing chronic bronchitis, a lowered risk of getting most common cancers (including breast cancer), as well as a lowered risk of getting cataracts. Additionally, fruits, like avocado, are high in protein, vitamin E, C, and beta-carotene and may be linked to stroke prevention.\nWhile food products comprising fruits and vegetables have been linked to health benefits in humans, such products are often difficult to prepare for sale in commerce. This is true because the quality of food products comprising fruit often deteriorates (e.g., browns, darkens, grows mold and/or loses flavor) due to enzymatic reactions within the food product, thereby resulting in a product that has a short shelf life and does not have an appealing look or taste after spending a limited period of time in conventional commercial channels.\nKnown techniques have been used to inhibit the deterioration of food products comprising fruits and/or vegetables. These known techniques include pasteurization of the fruit, high vacuum processing for removing oxygen, and chemically treating the fruit with sulfiting agents before making the food product. The above-described known techniques do not eliminate, for example, browning and darkening in food products comprising fruit, and such techniques have adverse effects on the flavor, aroma, texture and nutritional value of the fruits and vegetables treated, as well as the food products prepared therefrom.\nIt is of increasing interest to develop a ready-to-serve and ambient stable fruit-based composition that does not, for example, easily brown, darken and lose flavor and that has an extended shelf life at ambient temperature. This invention, therefore, is directed to a ready-to-serve and ambient stable fruit-based composition that has not been subjected to chemical treatment, high vacuum processing and temperatures over about 85° C. The ready-to-serve ambient stable fruit-based composition of this invention is low in fat, substantially free of starch, and suitable to be a beverage, filling, dip, sauce, spread, dressing or the like. Moreover, the ready-to-serve ambient stable fruit-based composition of this invention can be prepared from under ripe, ripened and/or over-ripened fruit and is ambient stable for at least about eight (8) weeks after opening. Such a fruit-based composition has substantially the same visual, texture, aroma and taste attributes of a fruit-based composition made on demand from freshly picked fruits."}
{"text": "The present invention relates generally to porous materials having high surface areas, and more particularly to such materials which are non-degradable, and methods for fabricating the same.\nMembranes have been typically used for filtration (microfiltration, ultrafiltration, nanofiltration), reverse osmosis (hyperfiltration), dialysis, pervaporation, and gas separation applications. See Scott, K. and R. Hughes, Industrial Membrane Separation Technology, 1996, London: Blackie Academic & Professional; Baker, R. W., Membrane technology and applications, McGraw-Hill professional engineering, 2000, New York: McGraw-Hill; and Cardew, P. T., M. S. Le, and Royal Society of Chemistry, Process Technology Group, Membrane processes: a technology guide, 1998, Cambridge: Royal Society of Chemistry.\nA solid membrane can be made of synthetic polymers, natural macromolecules, inorganic compounds, ceramic, or metallic materials. These membrane materials are generally fabricated through sintering, stretching, extrusion, phase inversion and etching, or casting. See Scott, K. and R. Hughes, Industrial Membrane Separation Technology, 1996, London: Blackie Academic & Professional; and Pinnau, I. and B. D. Freeman, Membrane formation and modification, ACS symposium series, 744, 2000, Washington, D.C.: American Chemical Society, Distributed by Oxford University Press.\nPorous membranes are advantageous in their low resistance to mass transfer of solutes in solution due to the increased permeation rate resulting from the pores. Therefore, porous membranes have been employed for separation of mixtures of proteins and macromolecules, salt concentration, sterilization, etc. They can also serve as 3-D matrices for chemical and biochemical mass exchange or reactions to take place, or for cells or other living organisms (e.g., bacteria, viruses, fungi) to grow. Therefore, they can be used as matrices in diagnostic systems, catalysis systems, culture systems, drug delivery systems, wound dressings, etc.\nU.S. Pat. No. 6,146,892 discloses a method for producing nanofibrillar matrices utilizing degradable polymers, such as for example, poly(L-lactic acid) (PLLA), poly(D,L-lactic acid-co-glycolic acid) (PLGA), and the like. The disclosed nanofibrillar matrices are highly porous and work well for various applications. However, the pore sizes are very small (on the order of 2 μm to 3 μm), which may in some instances render it more difficult for cells to enter. Further, small pore sizes may render it more difficult for material transport, especially materials which are particulate or contain particles. Yet further, the disclosed non-fibrillar structure (a platelet structure) was not as mechanically strong as may be desirable in some instances."}
{"text": "1. Field of the Invention\nThe present invention relates to a document data management system.\n2. Description of the Related Art\nIn an environment in which a document management system is used, a data of a document which is received via facsimile, a document which is scanned and transmitted by a multifunction peripheral or the like can be managed by storing the data to a folder of the document management system. A user carries out various kinds of operations with regard to the stored document data by making an access to a document management server from a client PC to display the prescribed folder. In this case, if a user carries out an operation which is the same as that he/she carried out in the past with regard to, for example, a stored facsimile document data, the user needs to select the operation again and execute the operation. That is, a user needs to repeat the same operation.\nHence, there has been proposed a document management system which can record information of an operation with regard to a document data as related information of a document, and select and acquire the information when other document data is operated and execute the operation.\nThere is a case in which a user combines a document data which is stored in a device of a multifunction peripheral or the like, and other document data which is stored in a document management system to be a single document, and carries out a collective operation with regard thereto. As a collective operation, for example, plural document data (file) are combined and converted into PDF, plural document data are made to be a binder, plural document data are summarized and transmitted by facsimile or an electronic mail and so on. However, in a case in which document data which are stored in different folders are summarized and subjected to a collective operation, it is necessary to copy the respective document data temporarily in a work folder (temporary work area), which takes time and labor. Also, in a case of carrying out the collective operation as described above with regard to a document data stored by using a sort server, it is conceivable that a combination of document data which are subjected to the collective operation differs for respective individual document data. In this case, it is necessary to carry out the operation by collecting document data of an object of the collective operation in accordance with respective document data.\nFor example, according to a technology which is described in Japanese Patent Laid-Open No. 2009-230221, although related information can be applied to a single document data, it is necessary to acquire and select the related information and execute the application at respective times thereof, and therefore, time and labor are taken."}
{"text": "One proposed power output apparatus is mounted on a vehicle that is equipped with an engine and a hydraulically-driven clutch to connect and disconnect the power transmission from the engine to drive wheels (see, for example, Japanese Patent Laid-Open Gazette No. 2003-74683). When the engine is separated from the power transmission to the drive wheels by means of the clutch and stops, this proposed power output apparatus performs a low pressure control to regulate the hydraulic pressure to a low hydraulic pressure level immediately before generation of an engagement force in the clutch. Upon satisfaction of engine restart conditions other than a vehicle start request, the power output apparatus performs feedback control to regulate the hydraulic pressure level immediately before generation of the engagement force in the clutch, based on the rotation speeds on both ends of the clutch. In response to a vehicle start request, the power output apparatus learns the regulated hydraulic pressure by feedback control as a hydraulic pressure for a next cycle of the low pressure control."}
{"text": "Many applications include instrumentation, which refers to an ability to: monitor or measure the application's level of performance, diagnose errors or record information regarding the execution of the application. The instrumentation can consist of application code for monitoring specific components in a system, for example, the code may output logged measurements associated with the performance of a system component to a log repository. An event log is a file (e.g., a table) with entries for events taking place in the execution of an application, such as transactions between a server and users of the application. The entry for each of the transactions (or other application activities) logged in the file is referred to as an “event”. An event log can include sequential entries for successive application activities as a series of nested events to indicate a hierarchy of events (e.g., an application activity that invokes another activity), with each entry describing a type, content, start time, end time, or other attribute of the application activities.\nEvent logs are important for understanding the activities of complex systems. However in complex client-server systems the number of transactions can be very large and it may be difficult to analyze the large volume of data in an event log for an execution of an application of such a system. For example, parallel processing of client-server transactions is often desired, however a failure to process application transactions in parallel (e.g., due to application misconfiguration or faulty application logic) can be difficult to detect from a simple examination of an event log for an execution of the application."}
{"text": "During surgery, a surgical drape is laid over a patient so as to cover a portion of the patient at or near a surgical field. A “surgical field”, which is also commonly referred to as a “surgical site”, is an operating field that is an isolated area where an invasive procedure or surgery is performed and that must be kept sterile by aseptic techniques. Surgical drapes are sterilized linens placed on the patient and around the surgical field in a manner that delineate sterile areas. Surgical appliances are placed on the surgical drape, such that the surgical appliances are within easy reach of the surgeon or the surgeon's assistant during surgery.\nThe surgical appliances used during surgery are typically scalpels for making incisions, retractors for holding open a portion of the body, forceps for holding organs and tissue, scissors for suturing and cutting, needle holders for holding needles while suturing tissue, tubing for allowing drainage or administration of fluids, staplers for closing the incisions, electrocautery devices for removing unwanted tissue or for sealing blood vessels, and other surgical appliances.\nIt is important that surgical appliances are contamination-free to reduce risk of patient infection. Contaminants (i.e., bioburden) that can cause infections include pathogens and microbial organisms, such as clostridia, streptococci, staphylococci, E. coli bacilli and other pathogens, bacterium and microbial organisms. Exogenous sources for these surgical field infectious pathogens, bacterium and microbial organisms include surgical personnel, the operating room environment, surgical instruments, and various materials brought to the surgical field. Such pathogens, bacterium and microbial organisms acquired by a patient while in the operating room can lead to serious post-operative, nosocomial (i.e., hospital-acquired) health complications, such as hepatitis, bronchitis, sepsis from intravenous sites, and even death. In the United States, approximately 780,000 of 30 million surgical procedures result in nosocomial surgical field infections. It has been estimated that nosocomial infections result in between 17,000 and 70,000 deaths annually in the United States.\nNosocomial patient infections may also result in medical malpractice tort liability for hospital personnel and/or the hospital facility, if medical personnel or the hospital fails to provide hygienic treatment and the infection is allowed to spread and cause further injury. In a 2005 study, it was estimated that the cost of jury awards for all medical malpractice cases in the United States was about 3.6 billion dollars. In 2006, the median medical malpractice award in the United States was $175,000.\nIn order to reduce risk of nosocomial patient infection during surgery, surgical appliances are routinely pre-sterilized. Especially in the case of re-usable surgical appliances, the surgical appliance is pre-soaked in a chemisteriliant solution and then hand scrubbed or subjected to ultrasonic cleaning. Pre-soaking, scrubbing and ultrasonic cleaning are performed so that any debris present on the surgical appliance cannot prevent direct contact between a sterilizing agent and microorganisms residing on the surgical appliance. In addition to pre-soaking, scrubbing and ultrasonic cleaning, pre-sterilization techniques also include chemical treatment, subjection to ionizing radiation, placement in a sterilizing chemical vapor or gas, and/or exposure to heat, as well as other pre-sterilization techniques. More specifically, chemical treatment may include use of chemicals with biocide capability, such as isopropyl alcohol, formaldehyde, bleach, tincture iodine, mercurochrome and other chemicals. Also, surgical appliances, especially surgical appliances having polymer components that cannot withstand elevated temperatures, can be exposed to a Cobalt-60 radiation source. The Cobalt-60 radiation source emits high-energy gamma ionizing radiation to kill microbial organisms. The radiation source can also be an electrical device that generates electron radiation in the form of an electron beam for killing microbial organisms. The chemical vapor or gas, which may be ethylene oxide gas, may also be used as a sterilizing agent. In the case of heat application, the surgical appliance is placed in an autoclave and subjected to moist heat in the form of pressurized steam or placed in an oven and subjected to dry heat. Dry heat may be applied at a predetermined elevated temperature, such as 320° F. (i.e., 160° C.), for a predetermined time duration, such as 60 minutes or wet heat may be applied at the same temperature of 320° F. (i.e., 160° C.), for a predetermined time duration, such as less than one minute. Use of the pre-sterilization techniques mentioned hereinabove reduces risk of patient infection when surgical appliances are delivered to and placed in the surgical field on the surgical drape.\nHowever, a surgical appliance can sometimes slip and fall from the surgical drape and land on the operating room floor. If this occurs, the pre-sterilized surgical appliance may become contaminated with infectious microbial organisms and may even break. In order to reduce the risk of patient infection and possible medical malpractice tort liability, the surgical appliance used during the surgical procedure must be replaced. In this case, the surgical procedure is interrupted, thereby causing more time to complete the surgical procedure. If sterilization of the surgical appliance is routinely performed inside the operating room, hospital personnel must repeat the previously mentioned sterilization procedures for the replacement surgical appliance. If the sterilized surgical appliance was delivered to the operating room contained in a package, another package containing a sterilized appliance must be obtained. Thus, re-sterilizing replacement surgical appliances in the operating room or obtaining another package containing a sterilized surgical appliance results in additional time to complete the surgical procedure, particularly if a sterile replacement surgical appliance is not immediately available. In addition, interruption of the surgical procedure can even pose a health risk to the patient. The health risk to the patient may arise because the surgeon will divert his eyes and attention away from the exacting surgery being performed in order to attend to replacing the surgical appliance. Therefore, it is important that surgical appliances are prevented from slipping from the surgical drape and falling to the operating room floor.\nVarious approaches have been attempted to address the issues mentioned hereinabove. For example, U.S. Pat. No. 4,944,311 titled “Surgical Instrument Retainer” and issued Jul. 31, 1990 in the names of Eldridge, Jr. et al. discloses a reusable, flexible surgical drape which is laid over a patient adjacent the surgical field and which retains surgical instruments thereon to facilitate access to the instruments. A plurality of magnets is embedded in the drape so as to retain magnetizable instruments placed on the drape by means of magnetic force. A non-megnetized portion is provided in the center of the drape for storage of non-magnetizable instruments. However, according to this patent, it appears that the non-magnetizable instruments are merely placed in the center of the drape rather than being securely attached to the center of the drape. Therefore, it appears possible that the non-magnetizable instruments can be inadvertently knocked or displaced from the center of the drape and fall to the operating room floor during the surgical procedure. Also, this patent appears directed to a surgical drape upon which surgical instruments are placed to facilitate access to the instruments and does not appear specifically directed to a surgical drape for preventing instruments from falling to the floor of the operating room.\nAnother approach is disclosed in U.S. Pat. No. 4,976,700 titled “Surgical Securing Tape” and issued Dec. 11, 1990 in the name of Dennis R. Tollini. The Tollini patent discloses a securing tape for securing to a patient's skin or to a support, a medical device such as tubing, a catheter, an intravenous needle, or the like. According to this patent, the securing tape includes an elongated tape having base portions and a central tab formed integrally therewith, The securing tape also includes pressure-sensitive tape on the base portions and on an exposed window of the tab. The securing tape further includes hook and pile fastener portions on opposite sides of the exposed adhesive on the tab and on the base portion facing the tape's exposed adhesive. However, it appears the Tollini patent is directed to securing tubes, catheters, intravenous needles, or the like and is not directed to securing larger surgical instruments of non-tubular shape, such as scalpels, retractors, forceps, scissors, staplers, and other larger, non-tubular instruments,\nYet another approach is disclosed in U.S. Pat. No. 5,315,985 titled “Endoscopic Instrumentation Kit And Package Therefor” and issued Mar. 21, 1994 in the names of Andre P. Decarie, et al. This patent discloses an endoscopic or laparoscopic instrumentation kit including at least one obturator and at least two sleeves forming a trocar assembly. An obturator is a removable plug used during insertion of tubular instruments. A trocar assembly is an assembly having a sharp-pointed instrument equipped with a cannula or tube and used to puncture the wall of a body cavity and withdraw fluid. The kit may also include a catheter, an endoscopic surgical instrument, tissue-gripping sleeve members and attachment devices for the trocar sleeves. The kit is packaged in a vacuum-formed enclosure having raised walls which correspond in size and shape to the instruments packaged therein for retaining and displaying the instruments. A method for utilizing the kit is also disclosed. Although Decarie, et al. disclose an endoscopic instrumentation kit and package therefor, the Decarie, et al. patent does not appear to disclose means for attaching the instruments to a surgical field.\nAlthough the prior art approaches recited hereinabove may disclose (1) a surgical drape including a plurality of magnets embedded in the drape to place magnetizable instruments placed on the drape; (2) a securing tape for securing a tubular medical device such as tubing, catheters, intravenous needles, or the like to a patient's skin or to a support; and (3) an endoscopic or laparoscopic instrumentation kit, the prior art recited hereinabove do not appear to disclose the invention described and claimed hereinbelow."}
{"text": "1. Field of the Invention\nThe present invention relates to a processing apparatus, a processing apparatus control method and a program.\n2. Description of the Related Art\nIn recent years, as information technology has advanced, technology has been being developed which is intended to improve input/output performances and reliabilities by connecting a processing apparatus and a storage apparatus with a plurality of input/output paths. See, for example, U.S. Pat. No. 6,526,521.\nIn the case where a processing apparatus and a storage apparatus are connected with a plurality of input/output paths, the processing apparatus selects one of the plurality of input/output paths to transmit a data input/output request to a storage apparatus. The selection of an input/output path is performed according to a round-robin scheme in which an input/output path is selected sequentially from the input/output paths, taking load dispersion into account."}
{"text": "Traditional methods of immunization are achieved by injection of a mixture of antibodies which immunoreact with an invading pathogen (i.e., passive immunization), or by vaccination, which stimulates the immune system to produce pathogen-specific antibodies. Since foreign antibodies are cleared by the recipient, passive immunity confers only temporary protection. Vaccination confers longer-lasting active immunity.\nIn order to be effective, vaccination must generate humoral and/or cell-mediated immunity which will prevent the development of disease upon subsequent exposure to the corresponding pathogen. The pertinent antigenic determinants must be presented to the immune system in a manner that mimics a natural infection. Conventional viral vaccines may consist of inactivated virulent strains, or live-attenuated strains (Old et al., Principles of Gene Manipulation: An Introduction to Genetic Engineering, Blackwell Scientific Publications, 4th edition, 1989). A general problem with using a vaccine consisting of a virus is that many viruses (such as hepatitis B virus) have not been adapted to grow to high titre in tissue culture and thus, cannot be produced in sufficient quantity (Id.). In addition, the use of inactivated viruses present a potential danger of vaccine-related disease resulting from replication-competent virus may remain in the inoculum. Outbreaks of foot-and-mouth disease in Europe have been attributed to this cause (Id.). On the other hand, attenuated virus strains have the potential to revert to virulent phenotype upon replication in the vaccinee. This problem has been reported to occur about once or twice in every million people who receive live polio vaccine (Id.). Moreover, encephalitis can occur following measles immunization with attenuated virus (Roit, I. M. Essential Immunology, Blackwell Scientific Publications, Sixth Ed., 1988). Another disadvantage of using attenuated strains is the difficulty and expense of maintaining appropriate cold storage facilities (Id.). A major disadvantage associated with the use of live virus vaccines is that persons with congenital or acquired immunodeficiency risk severe infections. Such persons include children in developing countries who are often immunodeficient because of malnutrition and/or infection with viruses or parasites (Id., Old et al., supra).\nAs a result of recent advances in molecular biology and peptide synthesis, it is possible to produce purified viral proteins or synthetic peptides for use in immunoprophylaxis (Murphy et al., “Immunization Against Viruses,” in Virology, Fields et al., Eds., Raven Press, New York, pp. 349-370, 1985). Purified antigens may be produced by synthesizing peptides which represent immunologically important domains of surface antigens of the pathogen. The synthetic peptide approach has been successfully used with an antigenic determinant of the foot and mouth disease virus (Id.). One problem with this approach is that the poor antigenicity of synthetic peptides has required the use of Freund's adjuvant to enhance the immune response in experimental animals (Id.). Since Freund's adjuvant cannot be used in humans, an effective adjuvant for human use must be developed (Id.). In addition, a single antigenic site may not be sufficient to induce resistance since large surface antigens usually contain several distinct immunological domains that elicit a protective humoral and/or cell-mediated response (Braciale et al., J. Exp. Med. 153:910-923 (1981); Wiley et al., Nature 289:373-378 (1981)). There may also be difficulties in stimulating an immunologic response to epitopes that are formed by noncontiguous parts of the linear protein molecule (Murphy, et al., supra). There is evidence that the majority of protein determinants are discontinuous and involve amino acid residues that are far apart in the primary amino acid sequence, but are brought into close juxtaposition by peptide folding (Roit, supra).\nThe alternative approach to preparing proteins for vaccines involves the use of cloned viral DNA inserted into a suitable vector to produce viral protein in prokaryotic or eukaryotic cells (Aldovini et al., The New Vaccines, Technology Review, pp. 24-31, January 1992). This approach, also, has several limitations. For example, one must devise suitable conditions for the optimal production of the recombinant protein of interest by the recombinant host cells. The protein product must be isolated and purified from the culture system, and obtained in sufficient quantities for use as a vaccine. Finally, it may be necessary to perform post-translational modifications of the purified protein (such as glycosylation and/or cleavage of a fusion protein).\nAn alternative to producing the recombinant antigen in vitro is to introduce nucleic acid sequences coding for the antigen into the cells of the vaccinee. In this way, the antigen is produced in vivo by the vaccinee's cells and provokes the immune response. Tang et al. (Nature 356:152-154 (1992)) have shown that it is possible produce an immune response to human growth hormone protein in mice by propelling gold microprojectiles coated with plasmids containing human growth hormone genomic sequences. The resultant variability in the production of antibody production was hypothesized to arise from the operation of the microprojectile device, or the coating of the DNA onto the microprojectiles.\nMore recently, Ulmer et al. (Science 259:1745-1749 (1993)) injected a plasmid carrying the gene for influenza A nucleoprotein into the quadriceps of mice. The mice produced nucleoprotein antibodies, indicating that the gene was expressed in murine cells. The mice also produced nucleoprotein-specific cytotoxic T lymphocytes which were effective in protecting the mice from a subsequent challenge with a heterologous strain of influenza A virus. Similarly, Wang et al. (Proc. Natl. Acad. Sci. USA 90:4156-4160 (1993)) observed that the intramuscular injection of a human immunodeficiency virus (HIV) type 1 envelope DNA construct in mice generated antigen-specific cellular and humoral immune responses. In addition, splenic lymphocytes derived from the inoculated mice demonstrated HIV-envelope-specific proliferative responses. Thus, direct inoculation of DNA coding for pathogenic antigens can provide an alternative to the use of viruses, proteins, or peptides.\nOne problem with using naked DNA for inoculation is the low efficiency of cellular uptake. For example, the protocol of Wang et al., supra, requires the injection of 100 μg of the DNA construct biweekly for a total of four inoculations. As described herein, the use of cationic lipids as a carrier for DNA constructs provides a more efficient means of genetic immunization. According to the present invention, genetic immunization can be achieved with as little as 5 μg of a DNA construct, which has been complexed with cationic lipid.\nLiposomes have been used as carriers of genetic information in the transfection of tissue culture cells. A fundamental problem of liposome-mediated transfection with liposomes comprising neutral or anionic lipids is that such liposomes do not generally fuse with the target cell surface. Instead, the liposomes are taken up phagocytically, and the polynucleotides are subsequently subjected to the degradative enzymes of the lysosomal compartment (Straubinger et al., Methods Enzymol. 101:512-527 (1983); Mannino et al., Biotechniques 6:682-690 (1988)). Another problem with conventional liposome technology is that the aqueous space of typical liposomes may be too small to accommodate large macromolecules such as DNA or RNA. As a result, typical liposomes have a low efficiency (Felgner, “Cationic Liposome-Mediated Transfection with Lipofectin™ Reagent,” in Gene Transfer and Expression Protocols Vol. 7, Murray, E. J. Ed., Humana Press, New Jersey, pp. 81-89 (1991)).\nLiposomes comprising cationic lipids interact spontaneously and rapidly with polyanions such as DNA and RNA, resulting in liposome/nucleic acid complexes that capture 100% of the polynucleotide (Felgner et al., Proc. Natl. Acad. Sci. U.S.A. 84:7413-7417 (1987); Felgner et al., Focus 11:21-25 (1989)). Moreover, the polycationic complexes are taken up by the anionic surface of tissue culture cells with an efficiency that is about ten to one hundred times greater than negatively charged or neutral liposomes (Felgner, “Cationic Liposome-Mediated Transfection with Lipofectin™ Reagent,” in Gene Transfer and Expression Protocols Vol. 7, Murray, E. J., Ed., Humana Press, New Jersey, pp. 81-89 (1991)). In addition, the polycationic complexes fuse with cell membranes, resulting in an intracellular delivery of polynucleotide that bypasses the degradative enzymes of the lysosomal compartment (Düzgünes et al., Biochemistry 28:9179-9184 (1989); Felgner et al., Nature 337:387-388 (1989)).\nVarious formulations of cationic lipids have been used to transfect cells in vitro (WO 91/17424; WO 91/16024; U.S. Pat. Nos. 4,897,355; 4,946,787; 5,049,386; and 5,208,036). Cationic lipids have also been used to introduce foreign polynucleotides into frog and rat cells in vivo (Holt et al., Neuron 4:203-214 (1990); Hazinski et al., Am. J. Respr. Cell. Mol. Biol. 4:206-209 (1991)). Therefore, cationic lipids may be used, generally, as pharmaceutical carriers to provide biologically active substances (for example, see WO 91/17424; WO 91/16024; and WO 93/03709). Thus, cationic liposomes can provide an efficient carrier for the introduction of foreign polynucleotides into host cells for genetic immunization.\nVarious cationic lipids are well-known in the prior art. One well-known cationic lipid is N-[1-(2,3-dioleoyloxy)propyl]-N,N,N-trimethylammonium chloride (DOTMA). The structure of DOTMA is: DOTMA, alone or in a 1:1 combination with dioleoylphosphatidylethanolamine (DOPE) can be formulated into liposomes using standard techniques. Felgner et al. (Proc. Natl. Acad. Sci. U.S.A. 84:7413-7417 (1987)) have shown that such liposomes provide efficient delivery of nucleic acids to cultured cells. A DOTMA:DOPE (1:1) formulation is sold under the name LIPOFECTIN™ (GIBCO/BRL: Life Technologies, Inc., Gaithersburg, Md.). Another commercially available cationic lipid is 1,2-bis(oleoyloxy)-3-3-(trimethylammonia)propane (DOTAP), which differs from DOTMA in that the oleoyl moieties are linked via ester bonds, not ether bonds, to the propylamine. DOTAP is believed to be more readily degraded by target cells.\nA related groups of known compounds differ from DOTMA and DOTAP in that one of the methyl groups of the trimethylammonium group is replaced by a hydroxyethyl group. Compounds of this type are similar to the Rosenthal Inhibitor of phospholipase A (Rosenthal et al., J. Biol. Chem. 235:2202-2206 (1960), which has stearoyl esters linked to the propylamine core. The dioleoyl analogs of the Rosenthal Inhibitor (RI) are commonly abbreviated as DORI-ether and DORI-ester, depending upon the linkage of the fatty acid moieties to the propylamine core. The hydroxy group can be used as a site for further functionalization, for example, by esterification to carboxyspermine.\nAnother class of known compounds has been described by Behr et al. (Proc. Natl. Acad. Sci. USA 86:6982-6986 (1989); EPO Publication 0 394 111), in which carboxyspermine has been conjugated to two types of lipids. The structure of 5-carboxylspermylglycine dioctadecylamide (DOGS) is: The structure of dipalmitoylphosphatidylethanolamine 5-carboxyspermylamide (DDPES) is: \nBoth DOGS and DPPES have been used to coat plasmids, forming a lipid aggregate complex that provides efficient transfection. The compounds are claimed to be more efficient and less toxic than DOTMA for transfection of certain cell lines. DOGS is available commercially as TRANSFECTAM™ (Promega, Madison, Wis.).\nA cationic cholesterol derivative (DC-Chol) has been synthesized and formulated into liposomes in combination with DOPE (Gao et al., Biochim. Biophys. Res. Comm. 179:280-285 (1991)). The structure of this compound is: Liposomes formulated with DC-Chol provide more efficient transfection and lower toxicity than DOTMA-containing liposomes for certain cell lines.\nLipopolylysine is formed by conjugating polylysine to DOPE. This compound has been reported to be especially effective for transfection in the presence of serum (Zhou et al., Biochim. Biophys. Res. Comm. 165:8-14 (1991)). Thus, lipopolylysine may be an effective carrier for immunization.\nIn addition, Gebeyhu et al. (co-pending U.S. application Ser. No. 07/937,508; filed Aug. 28, 1992) have developed novel cationic lipids according to the general formula: \nwherein R1 and R2 separately or together are C1-23 alkyl or alkyl or alkenyl, q is 1 to 6,\nZ1 and Z2 separately or together are H or unbranched alkyl C1-6 \nX1 is —(CH2)nBr, Cl, F or I n=0-6 or\nX2 is —(CH2)nNH2 n=0-6 or\nX3 is —NH—(CH2)m—NH2 m=2-6 or\nX4 is —NH—(CH2)3—NH—(CH2)4—NH2 or\nX5 is —NH—(CH2)3—NH—(CH2)4—NH(C2)3—NH2\nwhere p is 2-5, Y is H or other groups attached by amide or alkyl amino group or\nX9 is a polyamine, e.g., polylysine, polyarginine, polybrene, histone or protamine or\nX10 is a reporter molecule, e.g., biotin, folic acid or PPD, or\nX11 is a polysaccharide or substituted polysaccharide, or\nX12 is a protein or\nX13 is an antibody or\nX14 is an amine or halide reactive group or\nX15 is —(CH2)r—SH where r is 0-6 or\nX16 is —(CH2)s—S—S—(CH2)t—NH2 where s is 0-6 and t is 2-6.\nThese compounds are useful either alone, or in combination with other lipid aggregate-forming components (such as DOPE or cholesterol) for formulation into liposomes or other lipid aggregates. Such aggregates are cationic and able to complex with anionic macromolecules such as DNA or RNA."}
{"text": "This Nonprovisional application claims priority under 35 U.S.C. § 119(a) on Patent Application No(s). 092109286 filed in TAIWAN on Apr. 22, 2003, the entire contents of which are hereby incorporated by reference.\n(a). Field of the Invention\nThe invention relates to an illuminated logo device, more particularly, to an illuminated logo device for a display.\n(b). Description of the Prior Arts\nFollowing the rapid development of modern technology, the electronic products are becoming lighter, thinner, shorter, and smaller. Therefore, the urgent task of the manufacturers is: how to develop an electronic product having both artistic appearance and superior function, so that the product is most acceptable by public.\nBoth the desktop computer and notebook computer use displays for displaying data. Further, for the purpose of advertising, almost all manufacturers will design a logo on the casing of the display so as to distinguish it from other brands.\nA logo can have its special functions. For example, when the computer is placed in a dark location, an illuminated logo can be the best indicator showing the exact position of the computer. However, an additional light source inside the computer is required for illuminating the logo. Adding the additional light source will increase the cost and the size of these computers or other electronic products. Therefore, this situation is an urgent problem of the manufacturer that requires immediate solution."}
{"text": "The present invention relates to a new and distinctive soybean cultivar, designated S110192. All publications cited in this application are herein incorporated by reference.\nThere are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possesses the traits to meet the program goals. The goal is to combine in a single cultivar an improved combination of desirable traits from the parental germplasm. These important traits may include, but are not limited to, higher seed yield, resistance to diseases and insects, better stems and roots, tolerance to drought and heat, altered fatty acid profile, abiotic stress tolerance, improvements in compositional traits, and better agronomic quality.\nThese processes, which lead to the final step of marketing and distribution, can take from six to twelve years from the time the first cross is made. Therefore, development of new cultivars is a time-consuming process that requires precise forward planning, efficient use of resources, and a minimum of changes in direction.\nSoybean (Glycine max), is an important and valuable field crop. Thus, a continuing goal of soybean plant breeding is to develop stable, high yielding soybean cultivars that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties.\nThe soybean is the world's leading source of vegetable oil and protein meal. The oil extracted from soybeans is used for cooking oil, margarine, and salad dressings. Soybean oil is composed of saturated, monounsaturated, and polyunsaturated fatty acids. It has a typical composition of 11% palmitic, 4% stearic, 25% oleic, 50% linoleic, and 9% linolenic fatty acid content (“Economic Implications of Modified Soybean Traits Summary Report,” Iowa Soybean Promotion Board and American Soybean Association Special Report 92S (May 1990)). Changes in fatty acid composition for improved oxidative stability and nutrition are constantly sought after. Industrial uses of soybean oil, which is subjected to further processing, include ingredients for paints, plastics, fibers, detergents, cosmetics, lubricants, and biodiesel fuel. Soybean oil may be split, inter-esterified, sulfurized, epoxidized, polymerized, ethoxylated, or cleaved. Designing and producing soybean oil derivatives with improved functionality and improved oliochemistry is a rapidly growing field. The typical mixture of triglycerides is usually split and separated into pure fatty acids, which are then combined with petroleum-derived alcohols or acids, nitrogen, sulfonates, chlorine, or with fatty alcohols derived from fats and oils to produce the desired type of oil or fat.\nSoybeans are also used as a food source for both animals and humans. Soybeans are widely used as a source of protein for poultry, swine, and cattle feed. During processing of whole soybeans, the fibrous hull is removed and the oil is extracted. The remaining soybean meal is a combination of carbohydrates and approximately 50% protein.\nFor human consumption, soybean meal is made into soybean flour, which is processed to protein concentrates used for meat extenders or specialty pet foods. Production of edible protein ingredients from soybean offers a healthy, less expensive replacement for animal protein in meats, as well as dairy type products.\nThe foregoing examples of the related art and limitations related therewith are intended to be illustrative and not exclusive. Other limitations of the related art will become apparent to those of skill in the art upon a reading of the specification."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for making a hydrogenation catalyst having a base material coated with a catalytic metal using mechanical milling techniques. This method produces a hydrogenation catalyst that may be used for the remediation of a plurality of contaminated materials, including, but not limited to, polychlorinated biphenyls.\n2. Description of Related Art\nThe class of 209 aromatic chlorinated molecules resulting from the attachment of up to ten chlorine atoms to a biphenyl is collectively known as polychlorinated biphenyls (PCBs). Generally, PCBs are known to have the chemical structure C12H10-nCln and, in addition to other chlorinated synthetic aromatic compounds, are of great concern due to their toxicity and persistence in the environment. Among the properties of these synthetic colorless liquids are high chemical stability, low flammability, low thermal and electrical conductivity, and low solubility in water.\nTwenty-nine years following the 1976 Toxic Substances Control Act (TSCA) ban on their manufacture, PCBs remain a continued environmental threat. Their persistence owes to the very high chemical stability of these molecules. Prior to the TSCA ban, these favorable properties were exploited in a variety of applications including paint stabilizers, transformer oils, capacitors, printing inks, and pesticides.\nToxicity evidence was found adequate to justify TSCA regulation; however, the debate over the extent of PCB toxicity on organisms remains heated. PCBs are known to bioaccumulate and concentrate in fatty tissues. Studies suggest increased incidences of cancer with long-term PCB exposure. These studies are arguably inconclusive as they involve the simultaneous analysis of multiple congeners. Further complications arise from the potential for contamination of commercial mixtures with other more toxic chlorinated compounds such as polychlorinated dibenzodioxins (PCDDs) and polychlorinated dibenzofurans (PCDFs).\nUntil recently, only one option was available for the treatment of PCB-contaminated materials, incineration. This, however, may prove to be more detrimental to the environment than the PCBs themselves due to the potential for formation of PCDDs. PCDDs have been shown to exist at abnormally high levels in proximity to incinerators burning chlorine-contaminated materials. Cancer rates have also been “well correlated” to both dioxins and proximity to chlorine-contaminated waste burning incinerators. The temperature required to incinerate PCB-contaminated material may be up to 1200° C. Additionally, it is expensive to incinerate the PCB-contaminated material and to transport contaminated material.\nAn alternate approach that has been used to treat PCB-contaminated materials is bioremediation. This is accomplished by the reductive dechlorination of PCBs by anaerobic microorganisms. However, this usually results in an incomplete removal of PCBs. Furthermore, specific conditions are required for bioremediation and are not always seen in contaminated sites.\nSolvent extraction has also been used to treat PCB-contaminated materials. Organic solvents are used to extract contaminated matrix. However, extraction is only applied ex-situ and the PCBs are not actually destroyed but transported to another media that has to be further treated.\nBase-catalyzed decontamination has been used to treat PCB-contaminated material by adding NaHCO3 or NaOH to the media. However, this remediation process is only applied ex-situ due to operating temperatures.\nMetals have been used for the past 10 years for the remediation of halogenated solvents and other contaminants in the environment; however, zero-valent metals alone do not possess the activity required to dehalogenate PCBs.\nRecent literature reports the rapid and complete reductive dechlorination of PCBs using palladium coated iron (Pd/Fe) or magnesium (Pd/Mg) in aqueous medium. The use of these bimetallic particles for dechlorination relies on the reduction potentials of magnesium or iron coupled with the hydrogenation-type catalytic activity of palladium. The current techniques for preparing Pd/Fe material provides for the easy preparation of the material using pallameres. Although the Pd/Mg material has a greater thermodynamic driving force and forms a self-limiting oxide layer to prevent extensive air oxidation, this material cannot be easily prepared with pallamerse. Therefore, a process for making palladized zero-valent metals using cost-effective and efficient techniques would be valuable.\nCommunication theory is principally concerned with reducing the average size of particles in a sample of crystalline or metallic solid; however, it can also be used to understand mechanical alloying of particles. To accomplish either of these tasks, the most commonly used processes involve ball milling, vibrational milling, attrition, and roller milling.\nBall milling is a process in which a material is loaded into a canister partially filled with milling balls. The canister is then rotated at high speed on its major axis so that the balls are held by centripetal force to the inside wall until they reach the highest point inside the canister. Gravitational force then exceeds the upward force of the balls and they fall to the bottom of the canister where they impact other balls and the canister wall.\nIf some milling material is pinched between the participants in one of these collisions and the collision is of adequate energy, then the particles are fractured into smaller particles. It should also be noted that a certain critical speed of rotation is necessary for this process to occur. At lower speeds the balls will simply roll over one another and at extremely high speeds their upward force will exceed gravity at all points inside the canister and the balls will stick to the canister wall. The approximate critical speed of any ball mill is given below and is usually on the order of about 250 RPMs,Nc=7.05/(D1/2)                where Nc=critical speed and D=milling canister diameter        \nVibrational milling is a process similar to ball milling except that the milling vessel is vigorously shaken in a back and forth motion or in a back in forth motion in conjunction with a lateral motion that produces a “figure 8” path. This type of milling relies solely on the extremely high-energy collisions between rapidly moving milling balls rather than the collisions between the balls and the canister wall, as described for ball milling. Since vibrator mills can often shake canisters at a rate of approximately 1200 RPMs, often producing ball speeds of upwards of 5 m/s, vibrational milling commonly yields the desired reduction in particle size at a rate one order of magnitude faster than that of ball milling.\nTwo other less common types of milling are attrition milling and roller milling. These processes are not commonly seen in laboratory settings but are often seen in industrial work. Attrition milling relies on rapidly spinning paddles to stir the milling balls present in the milling vessel. The rate of size reduction observed is often similar to the rate of reduction observed for vibrator mills of similar size; however, due to the necessity of a cooling system this type of milling is often limited in its capabilities to systems that can be milled in liquid media.\nRoller milling is a process that relies on fracturing caused by stress induced in the system from the compression of materials between two rolling bars or cylinders. It is most often used for reduction of very coarse materials into less coarse materials that can later be reduced in size by other means. The previous two systems are not discussed quantitatively.\nFor all milling types (other than roller milling), the reduction of particle size relies on stresses induced in individual particles caused by collisions within the milling vessel. This process reduces the average particle size until equilibrium is reached, at which point no further size reduction is observed. This phenomenon can be explained if one considers crystal matrix impurities as the cause of fracture. As each fracture occurs at the point or plane of an impurity that disrupts the crystal structure, that particular discontinuity disappears, thus the average number of imperfections in the crystal structure of each particle in the sample reduces. Since the crystal structure of each particle in the sample becomes more perfect each time a fracture occurs, at some point the particles will exhibit a near perfect crystalline structure. At this time, the collisions within the mill will no longer be of adequate energy to cause fracture and the size equilibrium will be reached.\nAnother explanation for this phenomenon arises from the fact that it is more difficult to inflict upon smaller particles the necessary sheer required to cause fracture. This can be explained by the observation that the probability of a milling ball impacting a particle with the necessary directional velocity is reduced when the particle is smaller. Additionally, smaller particles have a higher surface activity and therefore, have a greater probability of being re-welded to form larger particles.\nThe equilibrium size of the particles in a milling batch has been the topic of much research and one of the more straightforward methods for calculating this size, based on a number of variables, follows and can be applied to all types of milling:\n      W    =                  mass  of  particles  with  surface  area            ≤      S                          W        i            ⁡              (        t        )              =                  ∑                  j          =          1                i            ⁢                        a          ij                ⁢                  ⅇ                                    -                              S                j                                      ⁢            t                                where            a      ij        =          [                                    0                                                              for                ⁢                                                                  ⁢                i                            <              j                                                                                          W                ⁡                                  (                  0                  )                                            -                                                ∑                                      k                    =                    1                                                        i                    -                    1                                                  ⁢                                  a                  ik                                                                                                        for                ⁢                                                                  ⁢                i                            =              j                                                                                          1                                                      S                    i                                    -                                      S                    j                                                              ⁢                                                ∑                                      k                    =                    j                                                        i                    -                    1                                                  ⁢                                                      S                    k                                    ⁢                                      b                    ik                                    ⁢                                      a                    kj                                                                                                                          for                ⁢                                                                  ⁢                i                            >              j                                          ]                  B      ij        =                                        ψ            ⁡                          (                                                x                  j                                /                                  x                  1                                            )                                δ                ⁢                              (                                          x                j                            /                              x                i                                      )                    β                    +                        (                      1            -            ψ                    )                ⁢                              (                                          x                j                            /                              x                1                                      )                    δ                ⁢                              (                                          x                j                            /                              x                i                                      )                    α                                S      j        =                  K        ⁡                  (                                    x              j                        /                          x              1                                )                    α      To use this equation, the desired mass of reduced size particles and the size of these particles is specified to yield the milling time required for achieving these parameters. Solving this series however, involves the use of many constants, which must be determined through calibration experiments using materials for which the constants are already known.\nIn other areas of research, such as mechanical activation of reactants in a milling vessel, researchers are focused more on the rate of particle size reduction as opposed to the average particle size of the end product. In general, the rate at which a mill reduces the average size of the particles being milled is a function of the probability of any particle being trapped between the participants in the above stated types of collisions when those collisions possess the energy necessary to fracture a particle. This focus has produced a number of functions, which have been supported empirically, to determine the rate of particle size reduction. A few of the more general examples are shown below for ball milling and vibrational milling.\nIn terms of the change in total surface area of materials                For Ball Milling        \n                    ⁢                  M        ⁢                              ⅆ            S                                ⅆ            t                              =              0.50        ⁢                  k          1                ⁢                  K          ϕ          ′                ⁢                  LJd                      -            1.7                          ⁢                  D          2.2                ⁢                  x                      β            +            1                          ⁢                              U            ⁡                          [                              1                -                                  1.1                  ⁢                                                                                    σ                        a                                            ⁡                                              (                                                                                                            k                              1                                                        ⁢                            xU                                                                                                              Y                                                              ρ                                B                                                                                      ⁢                            dD                                                                          )                                                                                    1                      /                      2                                                                                  ]                                1.5                                        ⁢    for                      ⁢          U      ≤      1                  M      ⁢                        ⅆ          S                          ⅆ          t                      =          1.1      ⁢              k        2            ⁢              K        ϕ        ′            ⁢              LJd                  -          2.2                    ⁢              D        2.7            ⁢              x                  β          +          1                    ⁢                        {                      1            -                          1.1              ⁢                                                                    σ                    a                                    ⁡                                      [                                                                                            (                                                      U                            +                                                          k                              1                                                        -                            1                                                    )                                                ⁢                        x                                                                                              Y                                                      ρ                            B                                                                          ⁢                        dD                                                              ]                                                                    1                  /                  2                                                              }                1.5                                ⁢    for                      ⁢          U      >      1                      For Vibrational Milling18        \n      M    ⁢                  ⅆ        S                    ⅆ        t              =      0.6    ⁢          k      1      ′        ⁢          K      ′        ⁢          J      2        ⁢          LD      2        ⁢          d              -        2              ⁢    αωθ    ⁢                  ⁢                  x                  β          +          1                    .                          ⁢                        U          ⁡                      [                          1              -                              230                ⁢                                                                            σ                      a                      ′                                        ⁡                                          (                                                                                                    k                            1                            ′                                                    ⁢                          xU                                                                                                      Y                                                          ρ                              B                                                                                ⁢                          d                          ⁢                                                                                                          ⁢                                                      α                            2                                                    ⁢                                                      ω                            2                                                                                              )                                                                            1                    /                    2                                                                        ]                          1.5            \nThese equations are applicable only when U≦1.\nwhere,U=volume of particles in mill/(volume of mill−volume of balls)                n=0.4d3x−3U        n=number of particles per ball, dimensionless        S=specific surface of particles, sq. cm./g.m         φ=ratio of mill speed to critical one, dimensionless        L=length of mill, cm.        J=fractional ball filling of mill, dimensionless        d=diameter of ball, cm.        D=diameter of mill, cm.        x,x′=particles sizes, cm.        Y=modules of elasticity of materials, g.f/sq. cm.        \nand all other terms are constants relating to the material being milled\nSince these equations are often difficult to use, empirical studies of milling variables may prove to be more useful.\nIn general, it can be seen that the rate of particle reduction in vibrational milling is much greater than in ball milling. Additionally, for all types of milling, rate increases with ball density and is greatest when the mill filling ratio (volume of material to be milled/volume of mill) is approximately 10-20% while the volume occupied by milling balls is approximately 40-60% of the total mill volume.\nIn the case where one wishes to mechanically alloy materials, the previously discussed theories apply, however several additional topics must also be considered. Mechanical alloying is a high-energy milling process for producing composite materials with an even distribution (though not homogeneous in the rigorous sense) of one material into another. By definition, at least one of the materials must be metallic to be considered an alloy; however, the topics discussed here can be applied to non-metallic materials as well. Two systems will be discussed; Malleable-Malleable and Brittle-Malleable.\nIn a malleable-malleable system such as the milling of two soft materials, like sodium and gold, ball-powder-ball collisions initially reduce the size of both materials until the average active surface area of each particle in the sample is large enough for re-welding to occur. When this critical surface area is achieved, in a particle, re-welding can occur between two similar particles or two dissimilar particles. If re-welding occurs between two similar particles, the net process results in no change of the material nature. If re-welding occurs between two dissimilar particles an alloy particle is created. This alloy particle can then undergo further fragmentation along alternative planes and subsequently be re-welded multiple times. The longer this process is allowed to take place the more dissolved one material becomes in the other.\nIn a brittle-malleable system such as palladium and magnesium, the more ductile magnesium is initially flattened while fragmentation of the more brittle palladium occurs. With further milling, the ductile material occludes brittle fragments leading to an individual particle composition of the starting mixture. Depending on the solubility of the brittle phase, continued milling may result in near chemical homogeneity.\nThese processes can be carried out using any type of mill however the lower energy route, ball milling, will require much longer milling times to produce a well-distributed alloy. This can be easily explained using theories already discussed. Simply put, ball milling produces relatively few collisions, most of which are relatively low in energy, while vibrational milling produces an abundance of very high-energy collisions of grinding material. Since high-energy collisions are necessary for alloying to occur vibrational milling produces the desired result with a greatly enhanced rate."}
{"text": "Penile compression clamps are well known and widely used to prevent urine leakage. In particular, penile clamps have been applied externally onto the penis for clamping the same. Such external penile clamps have been produced both as long-term reusable devices and also as short-term devices limited to one-week use. Typically, such devices have been employed for use in the inspection and treatment of diseases, wounds and other abnormal conditions of the bodies of humans and lower animals.\nFor example, Bard, Inc. produces the Cunningham clamp. The Cunningham clamp provides two pivoting arms connected by a single-hinge. Clamp arm surfaces that face the penis are padded for intended comfort. The Cunningham employs a ratcheted latch to clamp the two arms closer together and compress the penis. This clamp is reported to be a reusable, however, the large clamping forces often lead to pain, swelling and penile skin break down.\nThe Squeezer Klip™ disclosed in U.S. Pat. No. 6,463,932 provides a padded top and bottom arm in a pivoting configuration. Like the Cunningham, the Squeezer Klip™ is a hinged, reusable clamp. The dorsal and ventral arms include pressure-applying projections intended to contact the urethra, and preferentially apply pressure between the dorsal veins and arteries. The Squeezer Klip™ avoids direct compression of the neurovascular bundle. Patients are often incapable, however, of repeatedly applying the pressure-applying projections to the appropriate location to effect urethral closure. A screw mechanism clamps the penis using a fine compression adjustment via a threaded adjustment knob. The screw mechanism, however, requires a higher degree of manual dexterity to adjust compression, which can be difficult for older arthritic men.\nThe C3 clamp disclosed in U.S. Pat. No. 5,184,629 provides a clamp intended for use that is limited to approximately one week. The C3 is constructed from a co-extruded and thermoformed polyolefin sheet. The resultant component has two clam shelled halves connected by a hinge. The penis is placed through a portal between the two halves and the halves are folded over to compress the penis. The clam shells are held closed by a Velcro® strap. Force is localized on the urethra by presence of a raised bump on the bottom clam shell half. The C3 is available in only two fixed sizes. As there is great penile anatomical variation requiring multiple clamp sizes, a clamp frequently may be improperly selected. Further, the ability to operate the strapping system is often difficult for older, arthritic men. Such inconsistencies in sizing and the user dependent strapping system, however, make the C3 less reliable in its ability to control leakage.\nNone of the clamps above entirely eliminated urinary leakage.\nOther clamps have been disclosed in U.S. Pat. Nos. 6,609,522, 6,289,895 B1, 6,609,522 B2, and 6,234,174 B1. These patents, however, provide detailed straps, latches, hinges and mechanisms for “dialing” in required urethral pressure. Experience with the C3 and such clamps indicate that users may be confused by the added complexity of such components.\nWhile these clamps provide some advancement for controlling urinary dysfunction and protecting against bladder malfunction, improvements may yet be made to penile compression devices for males experiencing stress urinary incontinence. There is need for a penile compression device that provides optimum comfort and that is easy to apply and remove, while sufficiently preventing urinary leakage. Improvements may still be made to a penile compression device that provides an absorbent mechanism conveniently and comfortably attached to the compression device."}
{"text": "Tone mapping is a technique used in image processing and computer graphics to map a set of colors to another set so as to approximate the appearance of high dynamic range images in media with a more limited dynamic range. However, tone mapping may fail to reproduce the full range of light intensities present in natural scenes and may cause problematic contrast reduction from the scene values to the displayable range. Preserving the fine image details and color tones in the original scenes are important in many applications.\nRecently, Dynamic Backlight Control (DBC) in mobile Liquid Crystal Display (LCD) device has become one of the image display applications in applying the tone mapping for scaling up the image data while dimming the backlight so as to minimize the backlight power consumption. However, the image data loses its accuracy during the tone mapping especially in high brightness region.\nTherefore, the art would benefit greatly from image data processing methods that can avoid the loss of fine image details or contrast in dynamic backlight control applications."}
{"text": "1. Field of the Invention\nAspects of the present invention relate to a sulfur-containing mesoporous carbon, a method of manufacturing the same, and a fuel cell using the sulfur-containing mesoporous carbon, and more particularly, to a sulfur-containing mesoporous carbon having a high affinity to noble metal catalyst particles, a method of manufacturing the same, a supported catalyst using the sulfur-containing mesoporous carbon as a catalyst support, and a fuel cell using the supported catalyst.\n2. Description of the Related Art\nA catalyst plays an important role in fuel cell technology as the catalyst is contained in electrodes for the fuel cells to facilitate the electrochemical reaction therein; thus, many attempts have been made to increase the activity of the catalyst. Since the activity of the catalyst increases as the reaction surface area of the catalyst increases, the diameter of catalyst particles should be reduced to increase the reaction surface area of the catalyst and the catalyst should be uniformly distributed in the electrodes. As such, a catalyst support should also have a large surface area, and thus, much research has been dedicated to increasing the surface area of the catalyst support.\nIn addition to a large surface area, which may be obtained through high porosity, a catalyst support for a fuel cell should be electrically conductive so as to act as a path for the flow of electrons. Also, the catalyst support should support the catalyst particles such that the catalyst particles are highly dispersed, and the catalyst support should have a connected pore structure that facilitates the transfer and transport of fuels. Examples of such a catalyst support include amorphous microporous carbon powder, such as activated carbon or carbon black, structure-ordered carbon molecular sieve materials and the like. For example, see Korean Patent Laid-open Publication No. 2001-0001127.\nHowever, the micropores of such amorphous microporous carbon particles are poorly connected. Therefore, in a conventional direct methanol fuel cell (DMFC), a supported catalyst that is prepared using amorphous microporous carbon particles as a catalyst support exhibits a lower reactivity than unsupported catalyst in which a metal particle itself is used as a catalyst.\nHowever, if the metal particle itself is used as a catalyst, a large amount of the catalyst is required; thus, the manufacturing costs of the DMFC increase. Accordingly, the development of a supported catalyst that can improve the reactivity of a catalyst is urgently required.\nOne of the crucial problems of electrodes for fuel cells is an agglomeration of catalyst particles when driving the fuel cells for a long period or operating at a high temperature. The agglomeration of catalyst particles increases the size of the catalyst particles, which reduces the number of active sites and hence, the activity of an electrode reaction."}
{"text": "The introduction and temporary implantation through the skin, e.g., transcutaneously, percutaneously and/or subcutaneously, of biosensors has become very common in the treatment of patients inflicted with or suffering from any one of many different types of conditions. These implantable sensors include those monitoring a given parameter that indicates a certain bodily condition, e.g., a patient's glucose level, or the actual state of a treatment, e.g., monitoring the concentration of a drug dispensed to the patient or a body substance influenced by the drug.\nIn recent years, a variety of temporarily implantable sensors have been developed for a range of medical applications for detecting and/or quantifying specific agents, e.g., analytes, in a patient's body fluid such as blood or interstitial fluid. Such analyte sensors may be fully or partially implanted below the epidermis in a blood vessel or in the subcutaneous tissue of a patient for direct contact with blood or other extra-cellular fluid, such as interstitial fluid, wherein such sensors can be used to obtain periodic and/or continuous analyte readings over a period of time.\nCertain transcutaneous analyte sensors have an electrochemical configuration in which the implantable portion of these sensors includes exposed electrodes and chemistry that react with a target analyte. At an externally located proximal end of the sensor are exposed conductive contacts for electrical connection with a sensor control unit which is typically mountable on the skin of the patient. One common application of such analyte sensor systems is in the monitoring of glucose levels in diabetic patients. Such readings can be especially useful in monitoring and/or adjusting a treatment regimen which may include the regular and/or emergent administration of insulin to the patient.\nA sensor insertion device or kit is typically provided with such analyte monitoring systems for inserting the sensor into a patient. The insertion kit includes an introducer which typically has a sharp, rigid structure adapted to support the sensor during its transcutaneous insertion. Some introducers are in the form of needles having a slotted or hollow configuration in which a distal portion of the sensor is slidably carried to the desired implantation site, e.g., subcutaneous site, after which the insertion needle can be slidably withdrawn from the implanted sensor. Often, the insertion kit also includes an insertion gun for automatically or semi-automatically driving the introducer and attached sensor to within the skin. Implantation of a sensor with such an insertion device typically involves using the insertion gun to drive the introducer and pre-loaded sensor into the skin of the patient. The introducer is retracted into the insertion gun, leaving the sensor implanted within the patient.\nWhile such sensor insertion tools can greatly assist a user in effectively and efficiently implanting transcutaneous sensors, they are not without their drawbacks. As they are designed to be substantially automatic, the insertion guns tend to be mechanically complex, involving numerous static and moving parts. With the added complexity of such tools are significant costs in designing and fabricating them, contributing significantly to the overall cost of the sensor systems. In addition to the financial and manufacturing-related drawbacks, there are significant clinical consequences to using such transcutaneous sensor insertion tools.\nThe subcutaneous or other placement of such sensors, or any medical device, produces both short-term and longer-term biochemical and cellular responses which may lead to the development of a foreign body capsule around the implant. Consequently, this encapsulation may reduce the flux of an analyte to the sensor, i.e., may reduce the sensitivity or accuracy of the sensor function, often requiring numerous calibrations over the course of the sensor's implantation period. The extent of the immune response presented by implantable sensors, as well as the amount of pain and discomfort felt by the patient, are exacerbated by the size of the sensor introducer and/or the implantable portion of the sensor, often referred to as the sensor tail. With sensor introducers that carry the sensor within an interior or substantially interior space, there is naturally a limit on the extent to which the cross-sectional dimension of the introducer can be reduced.\nAccordingly, it would be highly desirable to provide a sensor introducer and associated sensor design, and/or a combined assembly, which do not require a separate insertion tool for their transcutaneous insertion, thereby minimizing the number of components involved, and reducing mechanical complexity and manufacturing costs. It would be additionally advantageous if the respective and combined dimensions and configurations of the introducer and sensor were further reduced to minimize the trauma, pain and immune response to sensor insertion/implantation."}
{"text": "The invention relates to a fuel cell system, in particular for a motor vehicle, having a fuel cell stack arranged in a housing, and having a measure for ventilating the housing by coupling a ventilation flow to another gas flow. With regard to the prior art, reference is made to DE 10 2008 016 579 A1.\nFor example, for use in a motor vehicle, it is known for a fuel cell stack to be arranged (to be “integrated”) in a substantially closed housing, wherein, aside from passage openings that are self-evidently required for media-conducting lines, it is also necessary for ventilation openings to be provided in a wall of the housing in order to permit ventilation of the interior of the housing. This is to be understood, in the present case, as being a “substantially closed” housing, wherein the present invention is also intended to encompass possible housing structures in the case of which no substantially closed side walls are provided, but with which a housing is realized which is substantially closed off at least in an upward direction and below the top side of which it is, for example, possible for hydrogen to accumulate. The hydrogen could, via possible leakage points, escape in small quantities from the media-conducting lines connected to the fuel cell stack.\nAn air flow, referred to in the present case as a ventilation flow, for the ventilation of the housing and discharge of possibly accumulating hydrogen from the housing is commonly introduced from the surroundings into the housing and discharged into the surroundings again. It is desirable for no dedicated (energy-consuming) delivery device to have to be provided for the delivery of the ventilation flow. In DE 10 2008 016 579 A1, it is proposed that the ventilation flow to be discharged from the housing be admixed to the cathode exhaust-gas flow of the fuel cell stack by way of a de Laval nozzle, utilizing the ejector pump effect. The cathode exhaust-gas flow or exhaust-gas flow of the fuel cell stack thus forms the delivery device for the ventilation flow.\nThe problem addressed by the present invention is to provide a further measure by which a ventilation flow through or out of a housing which contains a fuel cell stack can be discharged in an expedient manner.\nThe solution to the problem is, for a fuel cell system having a fuel cell stack arranged in a housing and having a measure for ventilating the housing by coupling the ventilation flow to another gas flow, characterized in that the ventilation flow is supplied to an air flow which is conducted through a radiator, in the form of an ambient-air heat exchanger, of the fuel cell system. Here, a radiator of the fuel cell system is to be understood generally to mean an ambient-air heat exchanger which is designed and provided for cooling at least sub-components of the fuel cell system, which may include not only the stack itself but also, for example, the power electronics of the fuel cell system.\nA fuel cell system according to the invention has a radiator in the form of an air-cooling fluid heat exchanger, wherein the cooling fluid is conducted in a suitable manner, for example through the fuel cell stack, and the waste heat thereof is dissipated at least partially via the radiator to the ambient air. The radiator exit air flowing through the radiator must in this case be discharged in a suitable manner, such that, for this purpose, an exit-air guide of a certain design is provided. In this case, the so-called ventilation flow is now also conducted away from the fuel cell system by way of the exit-air guide, such that the exit-air guide, in effect, performs a further function. Here, it may also be provided that a pressure gradient or the like, which at least partially contributes to the delivery of the ventilation flow, is generated by way of the ambient-air flow flowing through the radiator.\nAs used herein, the ventilation flow may be supplied or admixed, upstream of the radiator, to the air flow that is conducted through the radiator, because, upstream of the radiator, a directed flow is encountered which is not disrupted by lamellar structures or the like of the radiator and which, in effect, entrains the ventilation flow in accordance with the ejector pump principle. In this context, a ventilation opening of the housing can be positioned correspondingly, that is to say preferably directly adjacent to the flow path of the radiator inlet-air flow. In the conventional manner, it is then, or in this case, provided that, either for the radiator inlet-air flow or for the radiator exit-air flow, that is to say generally for the ambient-air flow to be conducted through the radiator, a suitable air delivery device in the form of a fan or the like is provided. In the case of a fuel cell system according to the invention installed in a motor vehicle, when the vehicle is moving, the “relative wind” also functions, in effect, as an air delivery device also for the ventilation flow of the housing that contains the fuel cell stack.\nSimilarly to the known prior art mentioned above, the ventilation flow of the stack housing may, however, also be supplied or admixed to an exhaust-gas flow of the fuel cell stack, for example by way of a Venturi nozzle, wherein, again, the ejector pump effect is utilized for the delivery of the ventilation flow. According to the invention, the exit-air flow of the radiator is then merged with the exhaust-gas flow, comprising the ventilation flow, of the fuel cell stack, such that multiple gas flows are, in accordance with the invention, conducted away from the fuel cell system together or jointly. Here, the exhaust-gas flow of the fuel cell system may be the cathode exhaust-gas flow and/or the anode exhaust-gas flow, which may by all means be merged, preferably downstream of the admixing of the ventilation flow, in a mixer.\nIn one aspect, a sensor for determining the hydrogen content in a gas flow downstream of the supply of the ventilation flow is provided, that is to say that, by way of a so-called hydrogen sensor of this type, the hydrogen concentration in the gas flow comprising this ventilation flow is measured. This gas flow may be an exhaust-gas flow as mentioned above or may be the radiator inlet-air flow or the radiator exit-air flow or a mixed gas flow comprising one or more of the flows. Thus, with the aid of an electronic monitoring unit, it can be checked whether an inadmissible amount of free hydrogen would pass with the mixed gas flow into the surroundings. It is possible, if appropriate, for suitable measures for this purpose, or for preventing this, to be introduced. If appropriate, with the aid of a sensor of the type and a suitable electronic monitoring unit, it is also possible, on the basis of the hydrogen content (=“sensor signal”) registered by the sensor, for a conclusion to be drawn regarding hydrogen possibly freely situated within the housing. If the gas flow contains exit air, already including hydrogen, of the fuel cell stack, it is specifically the case that the hydrogen concentration present therein is known, at least from a model calculation, to an electronic control unit for the control of the operation of the fuel cell stack, for which reason a hydrogen concentration higher than the known hydrogen concentration is indicative of additional hydrogen from the ventilation flow. It is basically thus possible, by way of a single sensor, to dispense with a (further) dedicated hydrogen sensor within the housing or for the ventilation flow alone.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of one or more preferred embodiments when considered in conjunction with the accompanying drawings.\nIn the figures, identical elements are denoted by the same references designations."}
{"text": "The present invention is concerned with a long-term feed for aquatic animals and especially for fish, shrimp and invertebrates, wherein the feed is stable for an extended time in fresh and sea water and can also be used as weekend and holiday feed for ornamental fish.\nLong-term feeds for ornamental fish, i.e., so-called feed blocks or feed stones, are known which have only a very small proportion of organic materials and mainly consist of calcium sulphate. DE 37 07 032 discloses a holiday feed which overcomes the long-term stabilization problems of the calcium sulphate feed blocks through the use of cement and crude rubber. A disadvantage thereof is the fact that the last-mentioned necessary additive materials are almost indigestable for aquatic animals and, in many countries, are not permitted according to the foodstuff laws.\nIt is an object of the present invention to provide an industrially producable, long-term feed for aquatic animals and especially for fish, shrimps and invertebrates, which feed is stable for an extended time and continuously provides ornamental fish with nutrients. The feed of the present invention can be used in fresh and sea water and is especially useful for warm and cold water ornamental fish without thereby making use of adjuvant materials which are not permitted in certain countries according to food-stuff laws but at the same time, however, does not impair the quality of aquarium water due to other adjuvant materials.\nSurprisingly, we have now found that, by means of the use of polysaccharides, for example of plant gums, the use of cement, calcium sulphate and/or latex can be avoided in the preparation of the feed. Additional properties which are substantially improved in comparison with the feeding agents described in DE 37 07 032 are also included in the present invention\nThus, according to the present invention, there is provided a long-term feed for aquatic animals, especially fish, shrimps and invertebrates in fresh and sea water. The feed is stable for an extended time., is insoluble, does not load the water chemistry in fresh and sea water and can also be used as a weekend or holiday feed for warm and cold water ornamental fish. The feed contains 1 to 50% by weight of a plant gum, preferred agar-agar. Preferred is a feed for ornamental fish which contains 1-10% plant gum, most preferred 2.9%. Typical plant gums which may be used are, for example, carrageenan, agar-agar, gum arabic, tragacanth, karaya gum, ghatti gum, carob bean meal, tara gum and/or xanthans. Preferably, the plant gum is agar-agar."}
{"text": "The present invention relates to a master cylinder configured to supply a liquid pressure to a braking cylinder of a vehicle.\nThis application is the U.S. national phase of International Application No. PCT/JP2013/058948 filed 27 Mar. 2013 which designated the U.S. and claims priority to Japanese Patent Application No. 2012-241198, filed Oct. 31, 2012, the contents of each of which are incorporated herein by reference."}
{"text": "It has been known that unnecessary electromagnetic field noises (hereinafter, simply referred to as noises) generated from electronic, information, communication, and industrial devices influence adversely, such as malfunction, on other electronic devices.\nFurther, it has been concerned that noises give some influences to a human body, and the influences given by noises have been studied.\nIn order to specify such invisible noise generation source, it is effective to visualize noises. Hitherto, devices to visualize noises have been proposed (refer to Patent Literatures 1 to 3).\nFor example, a device disclosed in Patent Literature 1 creates a diagram in which the respective electromagnetic field intensities (signal levels) of noises are displayed in a color on a two-dimensional surface. A device disclosed in Patent Literature 3 performs frequency analysis of noises on the surface of a measured object, and creates a contour map of noises of a specific frequency.\nA device disclosed in Patent Literature 2 determines the position of a sensor in a three-dimensional space around a measured object by parallax, and displays the respective electromagnetic field intensities (signal levels) of noises at the position in a color. Accordingly, a measurer can acquire information on the noises at the position of the sensor."}
{"text": "The present invention relates to internal combustion engines of reciprocating piston type and is particularly, though not exclusively concerned with engines of the general type disclosed in EP-A-0591153.\nThis prior document discloses an engine in which the or each piston is caused to move over at least a portion of the cycle at a rate which is such that the graph of its displacement against time differs from the sinusoidal shape which is inherently produced in conventional engines in which each piston is connected to a respective crank on a crankshaft by a respective connecting rod. In such a conventional engine attempts are made to match the combustion of the fuel/air mixture to the motion of the piston but the philosophy underlying the construction of the prior document is that the combustion is permitted to proceed in the optimum manner and the piston is caused to move in a manner which \"follows\" the combustion and is related to the nature and progress of the combustion process.\nMore specifically, the prior document discloses an engine in which the piston is caused to decelerate and thus to move more slowly than in a conventional engine at or around the point in the cycle at which ignition of the fuel/air mixture occurs and then to speed up again prior to reaching the top dead centre position (TDC). This is based on the recognition that in a conventional engine the piston is moving at substantially its maximum speed at the point at which ignition occurs and the compression ratio is altering at substantially its maximum rate and thus impedes the rate of propagation of the flame front through the fuel/air mixture and thus impairs the nature and completeness of the combustion process. However, slowing the piston down at around the ignition point means that the rate of increase in the pressure of the fuel/air mixture at the time propagation of the flame front commences is substantially less than is usual which results in the flame front propagating through the fuel/air mixture very much more rapidly than as usual.\nThe prior document also discloses that the piston is caused to reach its maximum acceleration and maximum speed at something between 0 and 40.degree. after TDC, instead of 90.degree. after TDC as in a conventional engine and thereafter to move more slowly than in a conventional engine in the latter portion of its working stroke prior to reaching the bottom dead centre position (BDC). This results in a decreased temperature of the exhaust gases and thus in reduced emissions of NOx and reduced erosion of the exhaust ports and valves.\nExtensive tests have been conducted on engines constructed in accordance with EP-A-0591153and these have shown that the engine does indeed have a substantially increased efficiency by comparison with conventional engines and also dramatically reduced emissions of unburnt hydrocarbons CO and NOx. Indeed, these tests have shown that the combustion process in the engines in accordance with the prior document proceeds in a manner which is fundamentally different to that in conventional engines, as evidenced by the fact that, for instance, the rate of pressure rise in the cylinder during combustion is about 6.5 bar per degree of rotation of the output shaft, as compared with about 2.5 bar in a conventional engine and that the combustion is complete within about 22.degree. rotation of the output shaft after TDC, as compared to about 60.degree. in a conventional engine.\nHowever, the engine disclosed in the prior document incorporates profiled cams cooperating with the pistons and not a conventional crankshaft and whilst such cams are wholly functional and technically satisfactory it would be preferable for the engine to incorporate a crankshaft of generally conventional type because mass manufacturing facilities for crankshafts are already available and the technology for manufacturing crankshaft type engines is more familiar and tried and tested than that for cam type engines."}
{"text": "Pain is the fifth vital signs besides body temperature, respiration, heart rate, blood pressure, which can serve as a warning when the body is hurt and cause a series of defensive responses of the body. However, excessive pain will cause damage to the body, and become one kind of unbearable torture to the body. Therefore, analgesia is an important task for healthcare practitioners. Chronic pain threatens the survival and quality of human life, and also causes a great burden on families and society. According to American Pain Society, it was reported that the prevalence of chronic pain in United States was estimated as 35.5%, including 105 million people. It costs over $ 100 billion, directly resulting in the consumption of healthcare spending and the loss of working hours.\nThe most commonly used analgesics include acetaminophen, non-steroidal anti-inflammatory drugs and opioid receptor agonists, such as tramadol and morphine. Non-steroidal anti-inflammatory drugs (NSAIDs) have weak analgesic efficacy, and may lead to gastrointestinal bleeding, renal failure and liver dysfunction as well as other side effects; and overdose of opioids will cause respiratory depression, and long-term use thereof will lead to side effects, such as constipation, abuse, dependence and addiction. In recent years, there are some other analgesics which are also marketed, for example antidepressants (such as Duloxetine, Lilly), anticonvulsants (such as Puri Galindo, Lyrica™, Pfizer) and selective cyclooxygenase-2 (COX-2) inhibitors (such as Parecoxib, Pfizer).\nA new mode for developing analgesic drugs is the mechanism-mediated development of new drug targets. For such mode, the mechanism of pain is mainly studied to find an important factor or target for regulating occurrence and development of pain, thereby developing a new medicament to intervene with the factor or target and achieve analgesia. Many clinical applications based on the results of such basic research are currently tested, however, few of them are successful.\nOne typical example is the development of anti-nerve growth factor (Nerve Growth Factor, NGF) antibodies. NGF is a member from neurotrophic factor family, and plays an important role in the survival and apoptosis of neuron during development. Basic research shows that NGF may promote pain, therefore blocking NGF may be an important method for treating pain. Therefore many anti-NGF antibodies were developed, including Tanezumab (Pfizer), SAR164877/REGN475 (Sanofi/Regeneron) and NJ-42160443/AMG403 (Johnson & Johnson/Amgen). These drugs show analgesic effects in animal models and clinical trials of phase I/II also show safety and analgesic effects, however, in the clinical trial of phase III in rheumatoid arthritis and ankylosing spondylitis patients for evaluating clinical efficacy and safety of analgesia of anti-NGF antibodies in combination with NSAID drugs, painless ischemic necrosis of caput femoris can be found in a small number of patients (Garber K Painkiller Novel Fate of the mAbs Hangs in Balance. Nat Biotechnol 2011, 173-174; Seidel M, Lane N, Control of Arthritis Pain with Anti-Nerve-Growth Factor, Risk and Benefit, Curr Rheumatol Rep, 2012, 583-585).\nTherefore, there is a huge market demand in the development of analgesics with new mechanism. Brain-derived neurotrophic factor BDNF is a neurotrophic factor discovered after nerve growth factor was discovered, molecular weight of which is 12.4 kDa. BNDF mainly distributes in central nervous system, and is also synthesized in peripheral nervous system and plays an important role in the regulation of survival, differentiation of neuron and synaptic plasticity and damage repair. Currently, there is an evidence to show that BDNF is not only an important factor in regulating development of nervous system and emotional disorders, but also an important mediator of pain.\nPrecursor for brain-derived neurotrophic factor (proBDNF) is synthesized in endoplasmic reticulum from BDNF gene through transcription and translation. The peptide chain is of 247 amino acids in length, and the theoretical molecular weight of the amino acid sequence is 27.8 KD. However, the molecular weight may be varied in a range of 32-36 kD due to different degree of glycosylation for the protein. In ProBDNF, the amino acid sequence at 1-18 positions is signal peptide sequence, and two fragments are produced during secretion process, wherein one fragment is a polypeptide fragment containing amino acids at 19-129 positions (that is, precursor domain), which is called precursor domain (proBDNF pro-domain), and another fragment is a fragment of amino acids at 130-247 positions (that is, mature domain), which will form BDNF with bioactivities upon processing.\nAt present, a lot of evidence shows that proBDNF is not only an intermediate of mature BDNF, but also can be used as a ligand to bind high-affinity receptor, p75 neurotrophin receptor (p75NTR) for exerting biological effects. Currently, the function of proB DNF-p75NTR signaling pathway in pain is unclear, however, there is a point-of-view that spinal sensitization mechanism of pain is similar to forming mechanism of Long-term potentiation (LTP)/long term depression (LTD) of hippocampus, while BDNF-TrkB and proBDNF-p75NTR signaling pathway are important signaling molecules for regulating LTP/LTD.\nIn 2012, experiments from Zhang, Yanling (Central South University) demonstrated that the expression of proBDNF in dorsal root ganglion was up-regulated after hind legs of a rat were cut; after an adenoviral vector of proBDNF was injected into hind legs of a rat, a great deal of proBDNF was expressed on nerve fibers by the adenovirus vector, and over-expression of proBDNF caused a significant decrease in withdrawal threshold of hind legs to mechanical stimulation and inflammation in plantar tissue; and withdrawal threshold of hind legs after the hind legs of a rat were cut can be increased by intraperitoneal administration of anti-proBDNF polyclonal antibody.\nCutting pain belongs to acute pain. For acute and chronic pain, both of outer peripheral sensitization and central sensitization can occur, and the difference between which is that: for outer peripheral sensitization, transmitters released by nerve fibers or axons in local tissue and intracellular signaling pathways are different, and pain transmitter and signaling pathways associated with dorsal root ganglion are also different, the latter of which lasts longer; in addition, for central sensitization, sustained activation of glial cells in spinal cord may also be the most important difference between the mechanisms of acute and chronic pain."}
{"text": "Use of imaging sensors in vehicle imaging systems is common and known. Examples of such known systems are described in U.S. Pat. Nos. 5,949,331; 5,670,935 and/or 5,550,677, which are hereby incorporated herein by reference in their entireties."}
{"text": "The following relates generally to wireless communication, and more specifically to persistent paging collision enhancement using dynamically switched higher order antennas in idle mode.\nWireless communications systems are widely deployed to provide various types of communication content such as voice, video, packet data, messaging, broadcast, and so on. These systems may be capable of supporting communication with multiple users by sharing the available system resources (e.g., time, frequency, and power). Examples of such multiple-access systems include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, and orthogonal frequency division multiple access (OFDMA) systems, (e.g., a Long Term Evolution (LTE) system, or a New Radio (NR) system). A wireless multiple-access communications system may include a number of base stations or access network nodes, each simultaneously supporting communication for multiple communication devices, which may be otherwise known as user equipment (UE).\nWhen a base station has information to transmit to a UE, the base station may transmit a paging message to prompt the UE to enter a connected state to receive the information. In some cases, a UE may have the capability of supporting multiple wireless services or subscriptions (e.g., radio access technologies (RATs), or subscriber identification modules (SIMs), etc.). Each of the multiple subscriptions supported by the UE may be associated with a different paging occasion or wakeup time during a reception time interval. In some instances, the UE may wake up for a paging occasion for each of the multiple subscriptions. In some such instances, conventional solutions for processing pages may be deficient."}
{"text": "In high speed magnetic tape handlers, magnetic tape is reeled back and forth between a supply reel and a take-up reel passing over a record/play-back head. The tape is drawn across the head by a capstan usually driven by a fast response low inertia DC motor. In such machines the tape is accelerated and decelerated at extremely high rates posing a very severe problem on means for maintaining uniform tension on the tape in the vicinity of the record/play-back head at all times. The inertia of the supply and take-up reels is such that, although they may be driven by powerful servo motors, they cannot provide the required uniform tension. Accordingly, it is usual to provide a substantial vacuum buffer on each side of the capstan and between it and the servo controlled reels. These vacuum buffers are typically provided by means of rectangular straight channels generally about 2 to 3 inches wide and as long as dictated by the response time of the reel servos. Vacuum ports at appropriate points, particularly at the bottom of the tank, couple a vacuum system to the buffer channels. Additional small buffers may be provided usually one on each side of the record/play-back head, to accommodate rapid acceleration of the tape.\nA typical high speed magnetic tape handler is constructed with the two reels (supply and take-up) mounted in the upper part of a cabinet enclosure with the capstan and head mounted below and the two vacuum tanks extending downward, centrally located and side-by-side. Illustrating the system with 101/2 inch diameter reels, a typical cabinet is of the order of 27 inches wide and 63 inches high. The result is a substantial area of unused space on each side of the vacuum tanks and below the reels. It is the object of the present invention to more efficiently use the cabinet area."}
{"text": "The instant invention relates to embryo transfer procedures. More particularly, the instant invention relates to a calculator for displaying synchronization of the ovulation cycles of a donor of at least one egg and a recipient of that egg during an egg transfer procedure.\nRecent advances in reproductive medicine and reproductive technology have provided novel avenues for the treatment of infertility and have made it possible to achieve pregnancy in women previously thought to be irreversibly sterile.\nThe development of human in vitro fertilization (IVF) has successfully overcome female sterility due to absent, obstructed or irreparable fallopian tubes and has resulted in the extension of this technique to any etiologic infertility factor which cannot be satisfactorily treated with conventional methods, provided the partners involved have normal gametes and the female partner has a normal uterus. Recently, successful transfer of donor embryos and donor oocytes fertilized in vitro to recipient endometria has extended the use of in vitro fertilization technology to conditions where female gametes are either absent or are not readily accessible to oocyte harvest or where there is gametic abnormality resulting in in vitro fertilization failure. Moreover, the availability of donor oocytes has allowed the use of such gametes for testing fertilization potential in couples where repetitive IVF trials have resulted in lack of fertilization of apparently normal oocytes by seemingly normal spermatozoa. Accordingly, oocyte donation programs can be utilized in the treatment of human infertility as well as diagnostically to assess spermoocyte interactions. Although, medically, oocyte donation is analogous to sperm donation, technically the relative inaccessibility of female gametes and the relative difficulty of synchronizing the ovulatory process in the donor with endometrial maturation in the recipient make the procedures quite different.\nWhile in vivo procedures have been successfully employed on a number of occasions, the transfer of in vitro fertilized donated oocytes to appropriately synchronized recipients appears to have wider clinical applications. In vitro procedures provide the added advantage of multiple oocyte collection in donors who are being stimulated with gonadotropines. Moreover, in vitro procedures do not require in vivo insemination of a donor with the inherent risks of undesired retained tubal or intrauterine pregnancy, and the possibility of transmitting semen born infectious diseases.\nSuccess of all donor oocyte procedures depends on appropriate embryo-endometrial synchronization. Critical to the achievement of pregnancy is the procurement of a viable embryo which can be transferred into a receptive endometrial milieu.\nIn order to maximize the possibility of a successful transfer, it is necessary to synchronize the development of the fertilized donor egg with the menstrual cycle of the recipient. While the subtleties of such synchronization may be readily understandable to physicians engaged in transplant research, the subtleties are rarely understood by recipients and donors and can be confusing to medical support personnel as well as physicians in the field. Since it is necessary to develop and communicate information regarding synchronization to recipients and their spouses as well as donors, there is a need for a procedure or device for conveying such information in a simple, readily understandable format which can eventually become uniform throughout the practice."}
{"text": "Although the epidemiology of headache disorders is only partly documented, taken together, headache disorders are extraordinarily common. It has been estimated that worldwide approximately 240 million people have migraine attacks each year. The National Headache Foundation states that more than 29.5 million Americans suffer from migraine headaches, with women being affected three times more often than men. In addition, in developed countries, tension type or “stress” headaches are estimated to affect two-thirds of all adult males and over 80% of adult females. Less well known is the prevalence of chronic daily headaches although the World Health Organization (WHO) estimates that one adult in 20 has a headache every or nearly every day. Trigeminal neuralgia is not a common disorder but the pain associated with trigeminal neuralgia attacks has been described as among the most severe known to mankind.\nNot only are headaches painful, but headache disorders can be disabling to afflicted individuals. Worldwide, according to the WHO, when analyzing all causes for “years lived with disability” migraine headaches were rated 19th on the list. Headache disorders may impose substantial hardships and burdens on the afflicted individuals including personal suffering, impaired quality of life and high financial cost. Repeated headache attacks, and often the constant fear of the next one, can damage an individual's family life, social life and their productivity at their place of employment. For example, it is estimated that social activity and work capacity are reduced in almost all migraine sufferers and in 60% of tension headache sufferers. Finally, the long-term effort of coping with a chronic headache disorder may also predispose an individual to other illnesses. For example, depression is three times more common in people with migraine or severe headaches than in healthy individuals.\nA wide range of headache types have been classified by the International Headache Society and among them are primary types including vascular, trigemino-autonomic and tension headaches and secondary types including headaches resulting from infection, trauma and non-vascular, intracranial disorders.\nVascular headaches refer to a group of headache conditions in which events at the interface between meningeal blood vessels and afferent nerve fibers are critical major components in the production of pain. Afferent nociceptive nerve fibers innervating meningeal blood vessels become activated in response to inflammatory and related events at the peri-vascular sheet causing a throbbing or pulsating type of pain. The most common type of vascular headache is migraine. Current pathophysiological evidence suggests that the trigemino-vascular system plays a pivotal role in the genesis of migraine headache. Migraine headache characteristics include unilateral (60%) or bilateral head pain, pain with a pulsating or throbbing quality, moderate to severe pain, pain associated with nausea or vomiting, sensitivity to light and sound, attacks that last four to 72 hours (sometimes longer) and visual disturbances or aura. Physical exertion often makes a migraine headache worse and women are more likely than men to have migraine headaches. Approximately one to two-fifths of migraine sufferers experience an aura, a sensory phenomena including visual disturbances that precede the onset of migraine headache. It is now believed that the aura is due to transient changes in the activity of specific nerve cells.\nAnother type of primary vascular headache is “cluster” headache, which is diagnosed by a well characterized clinical presentation. Although the syndrome is well defined from a clinical point of view, the causes are not well understood. The pathophysiology is believed to be associated with the trigemino-autonomic system. The periods during which cluster headache is experienced can last several weeks or months and then disappear completely for months or years leaving pain-free intervals between headache series. Cluster headache is characterized by frequent attacks of short lasting (15-180 minutes), severe, uniform, unilateral head pain associated with autonomic symptoms (e.g. lacrimation and nasal congestion). Pain can occur on the opposite side when a new series starts, pain may be localized behind the eye or in the eye region and may radiate to the forehead, temple, nose, cheek or upper gum on the affected side. Pain is generally extremely intense and severe and often described as a burning, boring, stabbing or piercing sensation. The headaches occur regularly, generally at the same time each day. Many individuals get one to four headaches per a day during a cluster period. Cluster headache is less common than migraine or tension headache and its cause is unknown. In contrast to migraine headaches, cluster headaches occur more in men than women and individuals suffering from these attacks may be very restless.\nTension headache, often referred to as “stress” headache, is a non-specific type of primary headache, which is of non-vascular origin and rarely is related to an organic disease. The pathophysiology of tension headaches is thought to involve the myo-facial system. For example, tension headache may be caused by the tightening of facial and neck muscles, clenching or grinding of teeth and/or poor posture. Tension headaches can be episodic and chronic in course, and typically are of mild to moderate intensity. Verbal descriptor used to characterize tension headache include pressing, dull aching and/or non-pulsating. Tension head pain is typically bilateral, and is not aggravated by physical activity.\nTrigeminal neuralgia, also called “tic duloreaux” is a condition that affects the trigeminal nerves and results in severe facial and head pain. Trigeminal neuralgia most commonly is diagnosed in patients over age 50, is slightly more common in women and has an incidence of approximately 4-5 per 100,000 persons. Trigeminal neuralgia is characterized by sudden severe, sharp facial pain, which usually starts without warning. The quick bursts of pain are described as “lightening bolt-like”, “machine gun-like” or “electric shock-like”. The pain is generally on one side of the face and is spasmodic, coming in short bursts lasting a few seconds which may repeat many times over the course of a day. Trigeminal neuralgia can involve one or more branches of the trigeminal nerve and the causes are not well characterized.\nCurrent Treatments\nThere are numerous treatment strategies for migraine and associated symptoms (e.g. nausea). However, to date, there is no single treatment strategy (including prevention or prophylaxis) that successfully alleviates migraine in a majority of patients. Additionally, treatment that has proven effective in one particular migraine sufferer may only be partially or intermittently effective. The current standard of care for migraine focuses on three major areas: 1) acute or abortive treatment; 2) treatment to relieve specific symptoms; and 3) preventive treatment.\nAbortive treatment is always indicated because of the disabling nature of migraine attacks. Sumitriptan and related 5-hydroxytryptamine (5-HT-1, serotonin) receptor agonists (triptans) are often considered the therapy of choice for migraine headache. To show optimal effectiveness, these agents generally have to be given early in the onset of pain. Serotonin receptor agonists are effective in up to 70% of patients and generally have few side effects when used sporadically. The number of patients benefiting from treatment with triptans may decrease to less than 50% during long-term therapy. Despite being effective, serotonin receptor agonists often only partially attenuate migraine headache. Rebound pain frequently occurs during the time period during which the drug levels are falling. Furthermore, side effects may arise including dizziness, heaviness or pressure on the chest and arms, shortness of breath, and sometimes chest pain which limit their clinical utility. Triptans are contra-indicated for patients with coronary artery disease. Other members of this class of drugs include, but are not limited to, sumatriptan, zolmitriptan, naratriptan, rizatriptan and elitriptan. Other classes of drugs that are used to treat migraine include the nonsteroidal anti-inflammatory drugs (NSAIDs), ergotamines and on occasion, neuroleptic drugs such as compazine. NSAIDS are as effective as triptans in alleviating migraine headache when given at the onset of mild migraine headache. Ergotamines, although commonly prescribed, are less effective than triptans and NSAIDS. Opioids such as codeine and butorphanol are not a first choice for the treatment of migraine because of their limited effectiveness, associated side effects including sedation and respiratory depression and the potential for dependence and abuse.\nPreventive treatment strategies are considered whenever migraine attacks have occurred several times in a month or are very severe and do not respond well to abortive medication. The following classes of drugs are used for preventing migraine: beta-blockers (e.g. propranolol, metoprolol, atenolol), calcium channel blockers, NSAIDs, antidepressants, anti-convulsant drugs (e.g. divalproex sodium, topiramate) and methysergide (no longer available in the United States). Another avenue for preventive treatment is educating the migraine patient to recognize and avoid migraine triggers which may help to reduce the frequency of attacks. Common migraine triggers include, but are not limited to, weather changes, bright lights, strong odors, stress and foods.\nIn a significant number of patients abortive and preventive treatments for migraine headache are often either ineffective, only partially effective or associated with significant side effects including hypotension, tiredness, increased weight, breathlessness, dizziness, heaviness or pressure on the chest and arms, shortness of breath, chest pain, nausea, muscle cramps, or peripheral vasoconstriction.\nTreatment strategies for cluster headache are classified as abortive or preventive. Abortive treatments are directed at stopping or reducing the severity of an attack, while preventive treatments are used to reduce the frequency and intensity of individual headache bouts. Abortive treatment strategies of cluster headache are quite successful when drugs can be injected. Drug classes used for injection include the 5HT1-agonists (e.g. sumatriptan) and the ergotamines. Alternatively, inhalation of 100% oxygen or an occipital nerve block have proven effective. However, all these treatment strategies require a visit to a doctor's office or to an emergency room.\nBecause of the short-lived nature of cluster headaches, preventive therapy is the cornerstone for individuals who have frequent attacks which severely affect their quality of life. Preventive therapy is initiated at the start of a cluster headache cycle and continues until the person is free of headaches for at least 2 weeks. The dosage of the preventive drug is then slowly tapered off which helps prevent relapsing headaches. Drug classes used preventively include beta-blockers, tricyclic antidepressants, anti-convulsants (e.g., divalproex sodium, topiramate), calcium channel blockers (e.g., verapamil), cyproheptadine, and NSAIDs (e.g. naproxen). Unlike drugs that ablate cluster headache (abortive drugs), most of the drugs used to prevent cluster headache have been developed for other clinical conditions and unfortunately, their effectiveness in prevention is limited.\nTreatment for tension headache usually consists of nonprescription painkillers such as aspirin, acetaminophen, NSAIDs or combinations of these agents with caffeine or sedating medications. When severe muscle contraction is present and/or the tension headache becomes chronic, more powerful prescription drugs may be needed to achieve relief. Tricyclic anti-depressants including amitriptyline HCl, doxepin HCl and nortriptyline HCl are commonly used. However these drugs have significant side effects including sedation, weight gain, dry mouth and constipation.\nThe first line treatment of trigeminal neuralgia is pharmacological in nature and is based on the use of antiepileptic agents including gabapentin, baclofen, clonazepam, lamotrigine, oxcarbazepine, toprimate and carbamazepine. About 50% of patients initially respond to treatment with a single agent and about 70% respond to treatment with two agents. However, a significant portion of patients (>50%) eventually becomes refractory to drug treatment and adding of a third agent or an analgesic drug (opioid or a non-steroidal anti-inflammatory agent) does not improve therapeutic success. Therapy with antiepileptic agents is also associated with side effects, most prominently dizziness, drowsiness, and ataxia. Many antiepileptic agents have the potential to cause rare but serious reactions.\nConsidering surgical interventions is the next appropriate step in patients who are refractory to pharmacological interventions. Surgical techniques include radio-frequency ablation of the trigeminal ganglion, micro-vascular decompression of the trigeminal root and gamma-knife radiation to the trigeminal root. The success rate of surgical techniques is initially quite high (80-90%) while the longer term success is closer to 50%. Specific side effects of these surgical interventions are sensory loss (numbness) and/or dysesthesia (e.g. analgesia dolorosa) in the distribution of the trigeminal nerve. Micro-vascular decompression in particular can be complicated by the occurrence of meningitis, cerebrospinal fluid leaks, or cranial nerve deficits. This procedure requires a craniotomy and the published mortality rate for this procedure is significant at 0.2-1.2%.\nIt is clear that migraine, cluster and tension headaches as well as trigeminal neuralgia can be debilitating to individuals and significantly impair their quality of life. To date, there does not appear to be a class of drugs or a treatment regimen that is effective for a majority of patients suffering from primary or secondary headaches or suffering from trigeminal neuralgia. Therefore, there is still a great need for novel and more effective therapies preventing or alleviating head pain of any origin."}
{"text": "This invention relates to packed bed assemblies of random-dumped packings used for mass and heat transfer operations between two fluids, typically operations such as gas absorption and desorption, distillation, liquid extraction, and the like.\nExtended-surface, random-oriented, packed beds for mass and heat transfer applications are widely used in industry, typically for gas-liquid and liquid-liquid contact. While most packed bed assemblies comprise vertical cylindrical columns employing countercurrent gas-liquid flow, referred to as packed towers, horizontal gas flow units are also known to the art, and are referred to as cross-flow contactors. To save energy and capital costs in either vertical or horizontal flow configurations, industry requires the highest flow capacities and lowest pressure drops for the packed bed contacting media. These objectives have been partially met in recent years primarily through the use of larger packing sizes.\nAt equal gas and liquid flows or loading rates in random-dumped beds, larger packing sizes of a given shape and design have relatively lower pressure drop and mass transfer performance than the smaller packing sizes. Because of their higher void space, both as individual packing elements, i.e, packing pieces, and as a packed bed, assemblies of the larger packing sizes also have higher limiting liquid and gas flow capacities (loading and flooding) than the smaller sizes of the same design. However, the gain in gas and liquid flow capacity for the larger packing sizes comes at the expense of lower mass transfer performance. The larger packing sizes have lower packing densities, i.e., lower number of packing elements per cubic foot and higher fractional void volumes, than do the smaller size packing elements of the same design. At the higher liquid loadings allowed by the lower packing density, the higher voidage of packed beds of the larger size packings results in increased gas axial back-mixing and loss of the desired countercurrent gas-liquid flow.\nVarious geometric packing shapes and configurations have been used in attempts to minimize the loss in transfer performance with increasing tower packing size. Nevertheless, for any given packing design, a significant portion of the gain in reducing tower diameter through the use of larger, higher flow capacity, packing sizes is offset by the deeper bed depth required to achieve the desired degree of mass transfer efficiency.\nMost random tower packings have shapes that provide anisotropic gas flow resistancexe2x80x94that is, they are relatively xe2x80x9copenxe2x80x9d for gas and liquid flow along one axis, and have a higher amount of planar or filamentary deflection surfaces along the transverse axis. Examples of packings with this property include cylindrical packings such as Raschig and Pall rings, spherical packings such as Jaeger Tri-Packs(copyright), as well as most plastic packings made by injection molding. To facilitate release from the mold, injection-molded packings necessarily have relatively high projected open area in the release direction, and typically, not in the transverse direction. Thus, these packings are highly anisotropic with respect to fluid flow resistance.\nWithin a packing element, gas and liquid flow will tend to take the path of least resistance, or through the xe2x80x9copenxe2x80x9d or mold release direction. While the random orientation of the elements in the bed is depended on to overcome this adverse property, gas flow macro-direction changes in the bed will be of the order of the packing size. In the smaller packing sizes, gas flow direction changes occur frequently with respect to the bed depth, and gas mixing is adequate. However, for the larger packing sizes the number of flow direction changes per foot of packing depth may be low enough to result in poor gas mixing. For example, a packed bed of 1-inch anisotropic packing elements would theoretically tend to have approximately 10-12 gas flow macro-direction changes per foot of packed depth. On the other hand, a 3xc2xd inch size packing would tend to have only 3 to 4 gas flow direction changes per foot of depth. This contributes to the loss in mass transfer performance with an increase in packing size in anisotropic packing elements. This characteristic has been tacitly acknowledged by the development of xe2x80x9cshallowxe2x80x9d packing shapes that have a relatively short axial depth, such as those disclosed in U.S. Pat. Nos. 3,957,931 and 6,007,915. There are additional parameters, such as the degree of liquid mixing and surface renewal frequency, whose contribution to transfer efficiency decreases with an increase in packing size for a given packing shape and style.\nThe use of sections or zones of different packing sizes has been taught in the prior art to resolve some specific packed bed problems. For example, because of the typically lower volumetric density of tower packing adjacent to the containing wall, the flow resistance in this area is less than in the center of the packing. Gas and liquid tend to channel along the wall, a phenomenon known as xe2x80x9cwall effectxe2x80x9d. Cameron and Bharga, in U.S. Pat. No. 5,679,290, teach the use of two different sizes of packing in order to overcome wall effect and to obtain uniform gas flow across the column diameter. In \"\"290, Cameron and Bharga use a plurality of a first packing size in an annulus adjacent to an upper part of the tower wall and a second plurality of a larger packing size in the core of the column in order to obtain uniform gas flow through the tower cross-sectional area.\nBecause flooding often occurs as a result of the restricted flow area adjacent to the packing support tray in countercurrent flow columns, it is known in the art to provide for higher flooding capacities by using a zone of a larger packing size on the support tray, and then a smaller packing size in the remainder of the column. It is also known in the art to use a zone or layer of smaller packing size on top of a larger column packing size in order to enhance initial liquid distribution. Various assemblies of zones and layers of packing elements are illustrated in the following Weber U.S. Pat. No. 2,055,162, Wible U.S. Pat. No. 2,271,671, Huber U.S. Pat. Nos. 3,285,587, 3,957,931, McKeown U.S. Pat. No. 4,002,705, Hoppe U.S. Pat. No. 4,333,894, Oshima U.S. Pat. No.5,242,626, Nagl U.S. Pat. No. 5,302,361, and Sunder U.S. Pat. No. 6,425,574. The zones of packing elements proposed in these and other prior art disclosures may be random or ordered, or combinations of one or more random or ordered zones or layers or packing elements. In none of the prior art of which I am aware, however, are two different packing sizes mixed with each other or completely co-mingled for use either as the sole packed bed contacting means or as the makeup of a bed, layer, or zone in a combination of beds, layers or zones.\nIt is an object of the invention to provide a method and apparatus for mass or heat transfer between two fluids in the form of a random-dumped packed bed having improved mass transfer performance combined with high limiting flow capacities.\nThis invention combines the supplemental surface area and mass and heat transfer capacity of a plurality of a first packing size with the lower pressure drop and increased gas and liquid flow capacity of a plurality of a suitable second larger packing size, by mixing the two sizes to obtain a substantially uniform mixture of the two disparate packing sizes. Mixing of the two packing sizes to a substantially uniformly dispersed mix may be done by any conventional means such as a tumbling drum, or controlled volumetric metering or a combination thereof. The mixed bed of this invention provides for the contribution of the supplemental mass and transfer area of the small packing size fraction while substantially retaining the higher limiting flow capacities of a bed of the larger packing size in which the smaller size packing elements are embedded.\nAdditional advantages accruing to the mixed-size packed bed of this invention include more uniform fluid flow, including possible alleviations of wall effect, and reduced axial back-mixing. A characteristic of a random-packed bed is that the size of the interstitial void spaces increases as the packing size increases. The interstitial voids are areas of zero flow resistance and zero active mass and heat transfer surface, except for liquid drop free fall therethrough. The interstitial voids between packing elements can also create low-resistance by-pass routes for gas and liquid, causing liquid channeling as well as gas axial back-mixing. These latter effects impair the desired countercurrent gas-liquid contact in a vertical packed column and result in lower contact efficiencies in a cross-flow contactor. Partially filling the interstitial inter-element voids of a packed bed of a larger size packing with smaller packing pieces replaces the empty void spaces with supplemental mass and heat transfer area, and provides for more uniform gas flow resistance and augmented liquid collection and re-distribution.\nBy studying shapes such as rings and spheres, it is possible to establish the general factors governing selection of the size mixture that will provide a uniformly dispersed and stable mixed bed in the method and apparatus of this invention. A mixture of packing sizes in which the volume of a smaller packing piece is approximately equal to, or greater than, the interstitial inter-element void space between the larger size packing elements provides mixed bed stability and is the preferred mixture. When the smaller packing element volume is larger than the average interstitial inter-element void space of the larger packing size bed, the final mixed bed volume is greater than the sum of the individual bed volumes of the two packing sizes prior to mixing. Because packings are sold on the basis of volume, this property is economically advantageous.\nCertain packing shapes or designs, particularly packing elements with external projections or ridges, tend to resist settling and may have unusually high interstitial or non-uniform inter-element void space. In such cases, attempts to mix in a packing size that is substantially smaller than the interstitial void space between packing elements may lead to separation, non-uniform distribution or difficulty in mixing. For example, attempts at mixing 1xc2xc-inch Jaeger Tri-Packs(copyright) with 3xc2xd-inch Jaeger Tri-Packs(copyright) gave beds where the smaller packing pieces fell through the interstitial voids of the larger packing and did not remain dispersed.\nWhile the size and shape of inter-element void spaces will vary with the design and size of the packing elements, void spaces are typically formed between more than two elements; that is, a grouping of elements. The number of the void spaces in a packed bed is therefore significantly less than the number of packing elements comprising the bed. Because the number of packing elements per unit volume increases rapidly with a decrease in nominal size, as is illustrated by the data of Table 1 for spherical and ring packings, relatively small volume percentages of a smaller packing size are required to substantially fill the void spaces of a random bed of a second larger size packing element.\nIn working with various combinations of packing elements to be used in my invention, the percent by volume in a bed and the percent by number of a given packing element in the bed can be calculated from a table such as Table 1, which can readily be made up for design configurations of packing elements other than those in Table 1. As an example, a 10% by volume mix of a nominal 1xc2xc-inch size Jaeger Tri-Packs(copyright) (1637 pieces per cubic foot) with a nominal 2-inch size Tri-Packs(copyright) (350 pieces per cubic foot) contains 162 1xc2xc-inch size packing elements and 315 2-inch packing elements, or a mix containing 34% by number of the smaller packing size. A 20% volumetric mix of the smaller 1xc2xc-inch size will contain 325 1xc2xc-inch pieces and 280 2-inch pieces, yielding a mix that contains 53.7% by number of the smaller packing size. Similarly, a mixture of only 10% by volume 1-inch Koch Flexirings(copyright) (1464 pieces/cubic foot) and 90% 1xc2xd-inch Koch Flexirings(copyright) (438 pieces/cubic foot) produces a packed bed containing 27% of the 1-inch rings by number.\nThere are a variety of packing shapes and geometries, including cylinders, spheres, saddles, etc., and the configuration and size of the interstitial inter-element voidage varies with the packing geometry and packing size. In addition to the well-known Jaeger Tri-Packs(copyright) and Koch Flexirings(copyright) mentioned above, see the variously configured packing elements disclosed in the following Ellis U.S. Pat. No. Ellis 3,957,931, Leva U.S. Pat. No.4,203,934, Hackenjos U.S. Pat. No. 4,203,935, Leva U.S. Pat. No. 4,316,863, Glen et al U.S. Pat. No. 4,554,114, Lang U.S. Pat. No. 4,724,593, Halbirt U.S. Pat. No. 5,543,088, Southam U.S. Pat. No. 5,670,095, Rukovena U.S. Pat. No. 6,007,915, Fan U.S. Pat. No. 6,182,950, Shojaie U.S. Pat. No. 6,371,452, and Sunder U.S. Pat. No. 6,425,574. My invention may be used with any of these or any other commercially available packing elements. This diversity of packing elements, however, suggests that preliminary testing of mixtures of a given packing design may be useful in order to suitable or optimum size mixtures for the application of this invention.\nIn the method and apparatus of this invention, the volume percent of the smaller packing size of the total packing volume is from 5 to 50%, and preferably from 5 to 30%. The appropriate volume and number mixture of sizes may be estimated from the piece count per unit volume data for each selected packing type and size combination by anyone skilled in the art of packed bed design. Alternatively, interstitial void volumes may be visually measured or estimated by mixing test quantities of the different packing sizes in a suitable container.\nPrecise calculations of the volumes of either the packing elements or the void spaces are difficult to make, particularly in view of the many different designs. Accordingly, as used herein, the term xe2x80x9csizexe2x80x9d is not an absolute, but a relative, term. The relative size of a packing element may be arbitrarily taken as equal to the volume occupied by a given number of randomly dumped packing elements divided by that number. For example, if 1000 packing elements A fill a cubic foot when randomly dumped, then the relative size of packing element A may be taken as equal to 0.001 cubic feet, or 1.73 cubic inches. Referring to Table 1, although they are both nominally designated as xe2x80x9c1-inchxe2x80x9d sizes, using the foregoing definition, the 1-inch Jaeger Tri-Packs(copyright) have a relative sized of 0.75 cubic inches as against 1.18 cubic inches for the Koch Flexrings(copyright).\nAs used herein, the term xe2x80x9cinterstitial inter-element void spacexe2x80x9d is also a relative termxe2x80x94that is, it represents an average of the volumes of the voids between contacting or nearly contacting packing elements. The average void volume between groups or assemblies of the packing elements is related to the relative size of the packing elements for a given configurationxe2x80x94that is, the larger the packing element size, the larger the void space between assemblies of the packing. Therefore, the relative interstitial inter-element voids may be taken as directly proportional to the relative size of the packing element as determined above. As used herein, when I refer to a packing element B which is of a size substantially equal to the average size of the interstitial void spaces of a random packed bed of a given type and size of packing element, I mean it has a size which, when distributed substantially evenly throughout a bed of the elements A, will occupy the void spaces in the bed of larger elements A while neither increasing the total volume of the bed nor falling through the void spaces.\nPreferred mixture percentages may also be based on the relative cost increase of the mixture compared to the base price of the larger packing size. For example, the prices of polypropylene Jaeger Tri-Packs(copyright), in the nominal 2-inch and 3xc2xd-inch sizes, are currently $14.24 and $9.36/cubic foot, respectively. From Table I, 5% by volume of the 2-inch size would have 18 elements, representing a cost of $0.72. The 95% of the 3xc2xd inch packing, or 46 pieces, would cost $8.97, giving a total cost for the mixture of $9.69. This is only $0.33 more per cubic foot than the $9.36/CF of the larger packing size, or a 3.1% increase, assuming volume additivity. However, tests show that, for this combination of packing sizes, the volumes are more than additive (the 2-inch polypropylene Jaeger Tri-Packs(copyright) packing is larger than the average interstitial voids of the bed of the 3xc2xd inch packing size) so that the percent cost increase would be even less than 3%. For this size mixture of polypropylene Jaeger Tri-Packs(copyright), there would be 18 of the 2-inch elements mixed with 46 of the 3xc2xd inch elements, giving a number percentage of 28% of the smaller packing size for a nominal 5% volumetric mix."}
{"text": "1. Field of the Invention\nThe present invention relates to a display apparatus having high resolution.\n2. Discussion of the Background\nWith the development of technology, the resolution of a liquid crystal display (LCD) has been gradually improved. Recently, a full high definition (FHD) LCD having a high resolution of 1920×1080 has been developed. However, since the LCD may have a hold type structure, motion blurring, in which an object is blurred when a dynamic image is displayed, may occur.\nIn order to prevent motion blurring, motion interpolation technology, which generates a new image frame having interpolated motion, and frame rate control technology, which adjusts the number of frames per second by inserting a new image frame between two sequentially input image frames, have been developed.\nHowever, a high resolution LCD employing motion interpolation technology has not yet been developed. Therefore, the display quality of a high resolution LCD may be degraded due to motion blurring."}
{"text": "1. Technical Field\nThe present invention relates to a system and technique for permitting a procedural or linear program to be accessed on an event-driven basis. More particularly, the present invention relates to a system and technique which permit programs written in procedural fourth-generation interpretive languages to be responsive to and be accessible by reactive or event-driven programs or input, such as a graphical user interface.\n2. Description of the Related Art\nA number of different varieties of programming languages exist today, each having its own attributes and characteristics. Among these languages are what are known as procedural fourth-generation interpretive languages. Programs written in such fourth-generation languages are not compiled, but rather are interpreted directly by a computer processor in the form in which they are written by a programmer, with the commands being executed by the processor based on the interpretation of the code written by the programmer. Such languages are typically easier in which to program than are languages which require compilation, and they are attractive because less programmer skill is required to write programs.\nPrograms written in procedural languages are executed linearly and are also known as command languages. The nature of procedural languages presents a barrier to some potential uses, because procedural programs typically cannot respond interactively to input from a user, but rather are restricted to being initiated through a command input by a user and then carrying out their programmed function to completion. While this type of performance is acceptable for certain functions or applications that computer users might want to have performed or carried out, such performance is not acceptable when users require object based event-driven programs which are reactive and capable of responding to any one of a variety of possible inputs/commands with the desired one of many possible functions or procedures. This requirement is the norm in user interface applications, and graphical user interfaces (GUIs) in particular.\nAttempts to make such fourth-generational procedural programs into event-driven applications have met with limited success. For example, one such technique has added a user interface module to the REXX (Restructured EXecutive eXecutor) fourth generational language which causes a user interface window to be displayed on a computer display device. However, this window merely permits a user to access a REXX procedural subroutine which continues until its processing is ended. Such single window modal user interfaces do not turn REXX programs or applications into the type of event driven programs or applications required by today's user interfaces, in which ease of use and a modeless design are prerequisites. For example, today's advanced GUIs permit users to customize or change them through working with the GUI. In contrast, the proposed modal REXX system would leave the user with the requirement of reprogramming the REXX subroutine itself in order to effect any change in the display of the window or in the function of the window.\nAccordingly, a need has arisen for a technique for converting or otherwise adapting fourth generation procedural language into event driven programs."}
{"text": "A sheet of the type of the present invention is disclosed by German Auslegeschrift 40 18 716. It is manufactured using a rubber-elastic mixture of special Shore A hardness of less than 85 which can be processed thermoplastically, and contains a thermoplastic polyurethane elastomer (TPU) and an ethylene/vinyl acetate copolymer (EVA). The vinyl acetate (VA) content of the EVA is taught to be between 78 and 95% by weight. In the rubber-elastic mixture there should be 50 to 99% by weight of TPU and 1 to 50% by weight of EVA, wherein the sum of the percentages is equal to 100%. The mixture is used to produce molded objects and sheets which are soft and elastic, and are readily processed (that is, without sticking), yet they also have a high strength, a good elongation at break, a high abrasion resistance, and show little swelling in fuels and lubricants.\nThe VA content of the EVA materials which are presently commercially available is of the order of 40 percent by weight. These materials can be calendered. If such an EVA is used in a thermally-processible mixture in accordance with the information of the German Auslegeschrift 40 18 716, the resulting mixture cannot readily be processed into sheets by calendering because the melt strength is too low. Moreover, calendering of this mixture results in a sheet with unacceptable sticking qualities.\nIt is therefore an object of the invention to develop an improved thermoplastic material having the desirable properties of the material of German Auslegeschrift 40 18 716, as well as having the capability of being processed into sheets, for example, by calendering."}
{"text": "Tools exist for producing computer-implemented graphical models. Using the tools, models can be created, edited, and executed within a computing system. These models can then be translated into a high-level language program code, such as C, C++, Ada, Java, or Javascript which may then be compiled into an object code or executable code. The object code or executable code can then be run on computing hardware independent of the graphical modeling tools.\nCurrent graphical modeling tools produce program codes that are generic in nature. For example, current tools may use program code that compiles or may be precompiled into object code or executable code which can be run on a wide variety of computing hardware. A single program code generated from a graphical model may compile into a single object code that may be run on, for example, an x086 Intel-type processor, an AMD-type processor, and a PowerPC-type processor. This generic single program code and the single object code the generic single program code compiles into may not be able to take advantage of specific performance enhancing features of these various processor types because the single program code was generated and compiled to run on all of these processor types.\nSingle program code compiled into the single object code may not achieve optimal or improvable performance on any computing hardware. Additionally, current graphical modeling tools may produce processor specific program code that may use processor-optimized function calls, inlined assembly code, or C/C++ language extensions in order to access the hardware specific optimization capabilities. For example, Simulink® and Real-Time Workshop® by The MathWorks, Inc. of Natick, Mass. (hereinafter “The MathWorks”) may provide target specific optimization capabilities for the C90 (ISO/IEC 9899:1990) and C99 (ISO/IEC 9899:1999) C-language standards as well as the GNU compiler tool chain. Additionally, target specific optimizations may be achieved with Simulink and Real-Time Workshop by creating a user-defined S-function block. However, these existing approaches may not be configurable, customizable or extensible through a documented API. Furthermore, their capability may be limited to predefined replacements and are not generally applied across all graphical modeling constructs."}
{"text": "1. Field of the Invention\nThe present invention relates to a folding top for a vehicle.\n2. Background Art\nA known folding top is borne by front and rear convertible top bows. The top bows are movably supported on a vehicle and extend transverse to the longitudinal direction of the vehicle. A housing having an interior light is mounted on a side of the front top bow facing the passenger compartment of the vehicle.\nDE 101 13 100 A1 describes a folding top having a flexible covering and a convertible top bow. The top bow holds the covering taut when the folding top is closed over a vehicle. The top bow has a rigid profile partly enclosing a hollow space. An opening in a boundary wall of the hollow space runs in the longitudinal direction of the top bow. An interior light is inserted into the opening to illuminate the passenger compartment of the vehicle."}
{"text": "The activities of a business or enterprise can be grouped into processes. Processes are business operations that are separated as desired and usually occur across business units. For example, the process of taking orders and turning those orders into revenue may be known as Order to Cash. The processes are comprised of sub-processes. For example, Order to Cash may be broken down into sub-processes such as Receive Order Inquiry, Provide Customer Quotation, Create Customer Outline Agreement, Create Sales Order, Schedule Production, Manufacture Product, Ship Product and Invoice Customer. Each sub-process may in turn be broken down into discrete activities such as providing customer number, entering that customer number, establishing pricing, determining a shipping date, etc.\nThe processes, sub-processes and activities operate, in part, by communicating information. For example, users may communicate through email. As another example, applications may communicate amongst themselves through electronic data interchange (“EDI”) and other similar services. Communication occurs horizontally, that is, among a process, sub-process and activities, as well as vertically, that is, between processes, sub-processes and activities.\nWhether communications occur horizontally or vertically, among applications or users, communications are increasingly asynchronous or message based. That is, enterprise communications were formerly primarily synchronous, or connection oriented, in which a connection is established with prior coordination between communication end points with data then being transmitted over the connection. Enterprise communications are now increasingly asynchronous, or connectionless, transmitting data without prior coordination between communication end points, such as through “event based” communications which use messages to move data instead of large files.\nAsynchronous or message based communications permit loosely coupled connections among and between systems because the end points do not have to be prepared to receive the data when the message is transmitted. Loosely coupled connections permit more flexibility in assembling processes. Flexibility in assembling processes is desirable in order to permit quick reaction to changing business conditions: if a particular sub-process or activity becomes unusable, the process can be reassembled with a new sub-process or activity. For example, if a Manufacture Product sub-process in the Order to Cash process at Widget Co. enterprise has a specific factory identified to manufacture the product and that factory has a fire or other disaster, making it unusable, Widget Co. will need to substitute a new factory. The ripple effect of that substitution among all of Widget Co.'s processes will change any number of parameters. A loosely coupled asynchronous connection among Widget Co.'s processes provides rapid substitution of the new factory for the old because the end points of communication to the new factory do not have to be predetermined before communications begin with the new factory. Thus, the flexibility of the asynchronous message based communication has permitted quick response to changing business conditions.\nDespite this flexibility, asynchronous or message based communications are problematic because of their loosely coupled nature. At any given time, precise information on the progress of the processes is difficult to obtain—messages may be in transit and not instantly locatable. For example, if a customer calls for the status of an order, an enterprise customer service representative may be able to determine nothing more than the fact that the order has been received and that the scheduled ship date is X. There is often no ability to drill down into the information levels and review the status in more granularity, such as the location of the good, the manufacturing status, etc., because the information required to review that status is in transit and unable to be reviewed.\nOf course, if the enterprise lacks the ability to access status information, business partners of the enterprise will similarly lack that ability. Thus, asynchronous communications may well increase inefficiency among business partners as well.\nThe difficulty in reporting caused by message based architecture also makes it difficult for the enterprise to measure the efficiency of its processes and their sub-process. Asynchronous messaging, with its indeterminate transmission of information, means a company may not be able to easily measure the interval between each sub-process, e.g. the time between Scheduling Production and the Manufacturing of a Product, and so easily measure the efficiency of their operations.\nFinally, asynchronous messaging may provide an enterprise with an ability to model and simulate processes. That is, since information flows can be readily estimated through enterprises with asynchronous messaging, and processes can be easily modeled from those flows, asynchronous messaging modeling provides the potential to model and simulate processes. That potential is not realized with present technology, however. Moreover, since as described above, enterprises lack information on the processes they have implemented, the enterprises are handicapped in their ability to modify those processes or plan new processes. A modeling and simulation tool, demonstrating processes, sub-processes and their activity or granular detail level would be extremely helpful, and would give the enterprise an opportunity to assemble, test, adjust, and simulate processes and their details.\nAccordingly, it is an object of the present invention to provide a tool for simulating message based architectures.\nIt is a further object of the present invention to provide monitoring capabilities for enterprise processes."}
{"text": "The present disclosure relates generally to a downhole nuclear magnetic resonance (NMR) apparatus, data processing, and interpretation methods for evaluating a characteristic of a region, and particularly for detecting and quantifying a gas-bearing earth formation in a subterranean region.\nNMR well logging is a technique used to investigate subterranean regions that may contain water, oil and/or gas reserves. The nuclei of chemical elements have a characteristic angular momentum (spin) and a magnetic moment, and by detecting and analyzing the reaction of the nuclei to applied magnetic fields, the characteristics of specific nuclei may be deduced. In the presence of an externally applied static magnetic field (B0), the nuclei spins become magnetized and align themselves parallel to the B0 field. By applying a radio frequency (RF) pulse train of a specific frequency to the magnetized nuclei, a pulsed RF magnetic field (B1) is generated that tips, or flips, the spins away from the direction of the B0 field. If the RF frequency (ω) substantially matches the condition for NMR (ω=γB0), where γ is the gyromagnetic ratio, then the first pulse (herein referred to as A pulse) reorients the magnetization to start precession and subsequent pulses (herein referred to as B pulses) generate spin-echo signals. A RF pulse sequence known as the CPMG (Carr-Purcell-Meiboom-Gill) sequence is typically used for well logging.\nAt the end of an A pulse, the spins are oriented transverse to the B0 field and precess around the direction of the B0 field at the Larmor frequency (ω0=γB0), and the transverse magnetization dephases with a transverse relaxation time constant (T2), also known as the spin-spin relaxation time. Repeated tipping and relaxation of the spins results in the NMR spin-echo signal, which may then be detected and analyzed for oil and/or gas field exploration.\nExisting methods use dual wait-time logs and the T1 contrast between gas and other formation fluids for gas detection and for gas saturation estimation. One such method is based on the assumption that water signals are fully polarized at both short and long wait time, TW, but gas signal is only partially polarized. So the difference between the two is contributed from gas only. However, if three phases coexist in the NMR detected sensitive volume, especially if slow relaxing water signal and light oil or oil-based mud filtrates are present in the formation, detection may be difficult or limited. Other methods for acquiring and processing multiple wait time data for T1 estimation employ improved log quality data using a summation of echoes approach, which is more useful if all of the partially polarized signal is gas. However, the certainty of the signal is reduced for discerning gas using this technique for wells drilled with OBM (oil based mud), and the subsequent invasion of the OBMF (oil based mud filtrate) into the sensitive volume, or for formations containing large-pore water or light oil. Thus skilled human interpretation is required in order to use existing art techniques. Accordingly, there is a need in the art for a robust NMR detection and analysis method that overcomes these drawbacks."}
{"text": "In the consumer products industry, individual products are typically placed in large containers or cartons that are shipped to retail establishments such as supermarkets. When a shelf needs to be stocked with additional product, an employee usually brings a full carton to the shelf, opens it, price marks each product inside, and places them one-by-one on the shelf, which is very time consuming.\nIn many instances, retail establishments prefer shipping containers that can be converted into display trays. Advantages of display trays include being able to place many products on a shelf in one motion, and being able to stack several open trays in an aisle if shelf space is not available. One such shipping/display container is a bottom tray filled with product that is shrink-wrapped with a web of thermoplastic film. The container can be opened by puncturing and unwrapping the film which is then discarded. However, most thermoplastic films are not biodegradable and therefore present an environmental concern.\nOne well-known type of shipping carton that can be converted into a display tray has printed instructions to circumferentially slice the carton's sidewalls along a dotted or dashed line with a knife, razor, or other sharp instrument. In opening the carton in this fashion, it is very easy to accidentally cut too deeply through the carton's walls and into the individual products inside the carton. If the product is solid or granular such as laundry detergent, the damage is generally confined to the outermost packages which, although messy and expensive, is not as catastrophic as slicing into individual containers that house a liquid such as dishwashing detergent. In addition, using a sharp instrument to cut open a carton poses a safety concern.\nOthers have identified the above-described problems associated with opening shipping cartons with a knife and have proposed various alternatives. For example, U.S. Pat. No. 3,850,363, which issued to Jacobs on Nov. 26, 1974, discloses a shipping carton having a rip cord circumferentially attached to the interior surface of the carton's side walls. Pulling the rip cord tears the carton's side walls and separates the carton into a top section and a bottom section. However, this general type of easy-open shipping carton has been found to be unacceptable for two major reasons; first, if the carton is made of corrugated paperboard and the rip cord is positioned perpendicular to the corrugated flutes, a very large pulling force is needed to split the carton walls; and second, it is nearly impossible to split the carton walls along a straight line which results in a mutilated bottom tray that is substantially weakened and unattractive and therefore generally unacceptable for display purposes.\nIn light of the above, a major object of the present invention is to provide a shipping container or carton that can be readily opened without having to use a knife or other sharp instrument.\nAnother object of the present invention is to provide a shipping container or carton that can be readily split along a pre-selected line into a top section and a bottom section, the latter being useful as an attractive display tray having a straight and smooth upper edge.\nYet another object of the present invention is to provide a high-speed manufacturing method of making blanks that are used in erecting shipping containers of the present invention.\nAnother object of the present invention is to provide blanks that can be erected into shipping containers by using conventional machinery.\nThese and other objects, advantages, and novel features of the present invention will become apparent to those skilled in the art from the following detailed description with reference to the drawings and appended claims."}
{"text": "Many catalytic processes are carried out in reactors which contain a series of separated catalytic beds. In such processes frequently mixing means, such as quench boxes, are disposed between beds. The purpose of the quench boxes is to provide rapid and efficient mixing of the fluid steams being processed in the reactor with a cooler fluid stream supplied from an external source thereby controlling the temperature of the process stream coming into contact with the catalyst in the next succeeding catalyst bed. As will be readily appreciated, the better the mixing, the better the temperature and reaction control and, hence, the better the overall reactor performance.\nExamples of interzone mixing devices can be found in U.S. Pat. No. 3,143,581, U.S. Pat. No. 3,480,407, U.S. Pat. No. 3,502,445, U.S. Pat. No. 3,705,016, U.S. Pat. No. 3,723,072, U.S. Pat. No. 3,880,961, U.S. Pat. No. 3,895,919, U.S. Pat. No. 4,235,847 and U.S. Pat. No. 4,836,989. Some of these devices are complicated. Some are prone to plugging. Some really need a relatively large space to provide the desirable degree of mixing.\nIn copending application Ser. No. 284,050 filed Dec. 14, 1988, now U.S. Pat. No. 4,960,571 there is described a quench assembly that provides a significant improvement over known quench boxes. Notwithstanding, the improvements achieved by that device, in today's competitive processing industry there is a need to constantly seek further improvement. It is, therefore, an object of the present invention to provide a device for thoroughly mixing in the space between two reaction zones, a stream of fluid supplied from an external source with the main process stream fluid.\nIt is also an object of the present invention to provide a device for efficiently mixing a relatively minor amount of fluid with a relatively major amount of fluid, especially in an interzone space. These and other objects are achieved by the novel device herein described and claimed."}
{"text": "A time-slotted channel hopping (“TSCH”) network can be defined by IEEE 802.15.4 and provide a communications network for resource providers (e.g., utility companies, home automation providers, industrial automation providers, or scientific and environmental application providers). A home area network may be used to communicate information between devices that consume resources (e.g., electricity) in a home or other premises and devices that monitor and/or manage the consumption of resources. Utility companies and other resource providers may use home area networks to monitor consumption of the resources by consumers.\nUtility companies and resource providers may operate on a time synchronized network, while a home area network may operate on a separate network that is not synchronized to the time synchronized network. Devices on the home area network may be idle for periods of time and communicate with the utility companies and resource providers intermittently."}
{"text": "(1) Field of the Invention.\nThe present invention relates to the crushing or digesting of used fluorescent tubes. More particularly the present invention relates to an improved rotatable crushing blade and the method of using such blade for the effective crushing of used fluorescent tubes into substantially uniform sized pieces of glass or pieces of glass within a uniform range of sizes prior to separating the glass and various potentially toxic materials contained within the interior of the fluorescent tubes from each other. More particularly, the present invention relates to an improved shape of a fluorescent tube fracturing blade and a range of speeds for operation of such blade which provides the desired results, i.e. the desired separation of the fractured glass particulates and the toxic powder and other substances contained originally within the fluorescent tubes.\n(2) Discussion of the Prior Art.\nFluorescent light tubes are formed from elongated, cylindrical or tubular glass receptacles which are charged with mercury or other conductive vapors. The inside surface of the tube is coated with a fluorescent coating of some form such as phosphorus itself or other phosphor powders such as antimony, beryllium, cadmium and strontium compounds plus in some cases, lead and the like. Mercury vapor as well as beryllium, strontium, lead and cadmium are well known as potentially toxic materials as are other phosphor powders with which the inside of the fluorescent tube may be coated. Older fluorescent tubes often use high concentrations of beryllium powders, but this has been in general superseded in more recent fluorescent tubes by cadmium-type powder.\nSince fluorescent tubes are, in general, bulky and unsatisfactory for disposal without treatment, it has become customary to crush them into small pieces by various means and then dispose of the fractured pieces. Merely fracturing the tube itself into small pieces for disposal, however, is not very satisfactory because of the potentially toxic nature of the dust and vapor originally confined inside the fluorescent tube. Such potentially toxic particulates, which occur mostly in the form of small dust particles plus mercury vapor and small drops or beads of mercury, can be quite detrimental if they escape to the environment.\nIn previous applications filed by the present Applicant, methods and means for crushing fluorescent tubes and then separating the fractured glass particulates from the potentially toxic materials contained within the original tube have been disclosed. Such methods and apparatus are based on an air separation effected between the glass particulates and the smaller toxic powders and vapors contained within the original fluorescent tubes. In particular, in accordance with such previous inventions, the fractured fluorescent tube materials are exposed to a rapidly moving body of gas such as stripping air, preferably passing countercurrently with such particulates. In this way the smaller, lighter dust and mercury vapor is carried away to a recovery system while the glass particulates are removed from the system for recycling to glass manufacturers and the like. In order to provide such separation, it is important that the glass particulates not be crushed too small so that they will have sufficient weight such that they will not be carried away by the stripping gas together with the small toxic powder particles.\nIn Applicant's U.S. application Ser. No. 458,177 filed Dec. 28, 1989 a flat, rapidly rotating blade was disclosed for fracturing fluorescent tubes thrust down a feed chute into the path of the rotating blade. The blade, upon striking the fluorescent tube, fractured such tube into small particles and these were then conveyed by gravity and other means through a system while separating toxic powder from the surfaces both by gravity processes and particularly by countercurrent air flow or gas stripping processes. As indicated above, it is important in crushing the glass that the pieces not be too small such that they may be carried away with countercurrent stripping air, nor should such pieces of glass be too large or they will not only not pack efficiently, but will also not have their surfaces exposed efficiently for transportation through the stripping air for stripping away of the toxic powder normally adhering to the original inside surface of the fluorescent tube. It is highly desirable, therefore, that the glass particles be of fairly uniform size and shape. It is also desirable, since glass is a very hard substance which rapidly wears away even fairly hard steel, for the rotating crushing blade to have a configuration and composition which will be durable and wear-resistant.\nIt has long been known to provide apparatus for generally breaking or crushing glassware, and particularly glass bottles and the like to reduce their bulk and more recently, it has also been widely suggested that fluorescent tubes be fractured into their constituent pieces for more convenient disposal as well as, in some cases, recovery of the components of the tubes. Among apparatus provided for crushing glass, and, of late, particularly fluorescent tubes, may be mentioned the following:\nU.S. Pat. No. 2,185,352 issued Jan. 2, 1940 to C. F. Peters discloses an early type of glass fracturing device for fracturing bottles. The fracturing device in the Peters patent is a pivoted hammer-type arrangement.\nU.S. Pat. No. 2,538,255 issued Jun. 26, 1951 to N. E. Johnson et al. discloses a remote control glass breaking machine in which bottles in particular are slid down a tube where they intercept a horizontally rotating motor driven fracturing blade comprised essentially of a rotating backing having slightly curved forward hammer surfaces which rapidly strike the bottle, pulverizing it by repeated blows. The hammer surface shown has a sharp edge on one side and is somewhat like a cleaver blade.\nU.S. Pat. No. 2,593,657 issued Apr. 22, 1952 to A. J. Coon et al. discloses a reciprocating-type crusher designed to crush fluorescent tubes and the like. The Coon et al. crusher involves reciprocating two opposed crushing surfaces relatively towards each other with a fluorescent tube between them, said reciprocating surfaces being moved by means of a rotating cam arrangement. There is no rotating fracture blade per se.\nU.S. Pat. No. 2,620,988 issued Dec. 9, 1952 to E. H. Telier discloses a fluorescent tube chopping device including a rotating fracture blade or \"star shaped breaking wheel\" which progressively breaks off the end of a fluorescent tube during the fracturing operation. Each arm of the fracture blade appears to have a substantially flat striking face.\nU.S. Pat. No. 2,628,036 issued Feb. 10, 1953 to J. B. Hall discloses a fluorescent lamp disposal arrangement including a rotating hammer-type arrangement for fracturing the fluorescent tubes. The rotating hammers, which are contained in a circular fracture chamber, are pivoted upon the outer circumference of a rotating disk.\nU.S. Pat. No. 2,866,604 issued Dec. 30, 1958 to J. B. Hall discloses a fluorescent tube disposal device including a rotating breaker arm.\nU.S. Pat. No. 3,333,777 issued Aug. 1, 1967 to C. W. Highfill et al. discloses a grinding mill which, although it is not designed for the fracturing of glass or fluorescent tubes, does disclose a series of flat blades arranged for progressively striking rocks and the like within a grinding chamber and in which fine dust is carried away.\nU.S. Pat. No. 3,353,756 issued Nov. 21, 1967 to D. J. Morgenson discloses a horizontally rotating hammer blade, into the path of which glassware such as bottles are dropped. The blade is rotated at a high speed such as 1700 rpm's so that a bottle dropped into the spinning blades may be struck as many as sixty times for each second the bottle remains in the contact zone. The blades of Morgenson are more like chopping blades than impact blades since they strike the bottles on their edge rather than on their flat side.\nU.S. Pat. No. 3,655,138 issued Apr. 11, 1972 to G. A. Luscombe discloses a multi-blade bottle or other glassware breaking device. The chopping blades comprise a series of fairly thin blades mounted in a stack, flat against each other to form a composite blade having in effect a number of knife blades sticking out from a central core at various points, which blades upon rotation of the central core strike anything passing by.\nU.S. Pat. No. 3,889,886 issued Jun. 17, 1975 to J. D. Spivey discloses a bottle breaking device including a series of edgewise rotating blades journaled on a transverse shaft. There are also a series of transverse stationary blades extending across the chopper and forming a grating between which the rotating blades pass during rotation. In effect, glass material has to be chopped within a size range which will slip down through the grating before it can pass beyond the chopping blades.\nU.S. Pat. No. 3,913,849 issued Oct. 21, 1975 to I. M. Atanasoff et al. discloses a fluorescent tube digester or breaker in which the breaking of the fluorescent tubes is accomplished by a double bladed knife rotating horizontally next to the tube inlet.\nU.S. Pat. No. 4,545,540 issued Oct. 8, 1985 to A. Nakamura discloses a fluorescent tube breaking device in which the tubes are broken by rotating fracture blades which rotate edgewise on a shaft in cooperation with a number of other blades and catch the fluorescent tube against a series of stationary blades mounted effectively between the rotating blades. The arrangement is substantially similar to the Spivey patent cited above.\nU.S. Pat. No. 4,579,287 issued Apr. 1, 1986 to W. E. Brown discloses a fluorescent tube fracturing device including a pair or plurality of flailing chains attached to a rotating disk. Both the disk and the chains rotate horizontally and intercept a vertically inserted fluorescent tube.\nU.S. Pat. No. 4,607,798 issued Aug. 26, 1986 to K. F. Odlin discloses a lamp crushing apparatus having a special allegedly non-jamming crushing blade which takes the form of a plurality of surface compartments in a drum section into which compartments or pockets formed between vanes at the surface of the drum the lamps are inserted, crushing said lamps into more or less equal, discrete quantities of fractured glass.\nU.S. Pat. No. 4,655,404 issued Apr. 7, 1987 to J. W. Deklerow discloses a fluorescent tub chopping apparatus in which the tubes are inserted into the path of a rotating chopping means formed from a pair of flails comprised of rectangularly shaped bars secured to a central rotating plate by hooks. The tube to be fractured is inserted past the horizontally rotating flails.\nU.S. Pat. No. 4,786,000 issued Nov. 22, 1988 to E. P. Weil et al. discloses a bottle breaking apparatus comprising a pneumatic ram with a central punch and a following plate which fractures the bottle against a V-shaped supporting wall. It is said the device does not form as many small particles of glass which may be difficult to recycle. The plate is deliberately not completely advanced against the bottle to avoid crushing said bottle into small pieces.\nU.S. Pat. No. 4,819,883 issued Apr. 11, 1989 to E. P. Weil et al. uses the same pneumatic punch and following plate to crush glass bottles as the prior Weil patent and in addition provides an angular support wall in the rear that does not support the bottle in the center and increases its shattering.\nWhile the above devices have generally been effective to fracture fluorescent tubes as well as other glass materials and other compositions of materials into small pieces, such blades generally have not been effective to provide a uniform fracturing of the tubes."}
{"text": "The present invention relates generally to methods and apparatuses for performing an infrared treatment, and more particularly, to methods and apparatuses for performing an infrared treatment on extruded plastic and elastomeric products in an infrared oven.\nIt is known in the prior art to heat plastic or elastomeric product that comes out of the extruder to cross-link the plastic or elastomeric product to obtain a desired effect that changes the capabilities or properties of the product (e.g., increases the strength of the product, changes the product to a solid, etc.). Furthermore, it is known to perform this heating of plastic or elastomeric product using infrared radiation and to use a wavelength of the radiation that penetrates inside the wall of the product in such a way that the heating takes place at the same time in all the depths of the product.\nOne prior art method is disclosed in patent publication GB 2283 489. In this patent publication, the material is cross-linked by using infrared radiation in such a way that the temperature obtained by the radiation corresponds to the wavelength demanded by the cross-linkage reaction.\nA similar prior art method is presented in the patent publication FI 109706. According to this method, an additional material or additive that is used for modifying either physically or chemically the properties of the plastic material is decomposed utilizing infrared radiation, wherein the wavelength is selected so that the radiation penetrates through the plastic material itself as efficiently as possible, but is also absorbed by the additional material, thereby heating and decomposing the additional material.\nThe greatest disadvantage of the prior art methods is that the infrared radiation inevitably consists of a distribution of different wavelengths. It is also inevitable that a part of this wavelength distribution follows approximately part of the curve of Gauss, and the wavelength distribution has rays of long wavelengths that do not penetrate into the material. Rather, the long wavelengths get absorbed by the surface of the product, causing inconvenient overheating of the product. Overheating causes the surface of the product to become oxidized or to react in some other unwanted way.\nAttempts have been made to solve this problem. For example, U.S. Pat. No. 6,106,761 (“the '761 patent”) addresses these issues by eliminating the infrared rays that correspond to absorption peaks of the material to be heated in order to minimize overheating of a surface of the material. The '761 patent notes that eliminating these rays may be accomplished by filtering out the unwanted rays. The filtering process disclosed in the '761 patent is very difficult to undertake because, when filtering out certain wavelengths, the filter itself gets overheated and becomes a source of infrared energy that sends the same filtered wavelength to the material, thus overheating the material.\nOne solution that attempts to avoid the overheating caused by filtering is cooling of the surface of the material during the infrared treatment. This can be done, for example, by blowing cool gas, like air, on the material. The greatest disadvantage of this method is that the air also cools the infrared lamps and reduces the capacity of the lamps. Another disadvantage is that dirt and other debris splashes from the material to the lamps and, thus, the lamps get dirty, which again reduces the capacity of the lamps.\nIt is very important for the irradiated material to be heated uniformly across an entire cross-section of the material. A method and apparatus for performing an infrared treatment, for example, on plastic and elastomeric products, that overcomes all of the previous obstacles and that uniformly heats the product is therefore desired."}
{"text": "In the related art, an image forming apparatus can perform printing using a color material that may be decolorized (decolorable material) and a color material that cannot be decolorized (non-decolorable material).\nSuch an image forming apparatus typically receives a print job from a printing processing device such as a personal computer, for example, and forms an image based on the received print job. The print job is usually generated by a printer driver executed in the printing processing device."}
{"text": "Multi-processor systems employ two or more computer processors that can communicate with each other, such as over a bus or a general interconnect network. In such systems, each processor may have its own memory cache (or cache store) that is separate from the main system memory that the individual processors can access. Cache memory connected to each processor of the computer system can often enable fast access to data. Caches are useful because they tend to reduce latency associated with accessing data on cache hits, and they work to reduce the number of requests to system memory. In particular, a write-back cache enables a processor to write changes to data in the cache without simultaneously updating the contents of memory. Modified data can be written back to memory at a later time.\nCoherency protocols have been developed to ensure that whenever a processor reads a memory location, the processor receives the correct or true data. Additionally, coherency protocols help ensure that the system state remains deterministic by providing rules to enable only one processor to modify any part of the data at any one time. If proper coherency protocols are not implemented, however, inconsistent copies of data can be generated.\nThere are two main types of cache coherency protocols, namely, a directory-based coherency protocol and a broadcast-based coherency protocol. A directory-based coherency protocol associates tags with each memory line. The tags can contain state information that indicates the ownership or usage of the memory line. The state information provides a means to track how a memory line is shared. Examples of the usage information can be whether the memory line is cached exclusively in a particular processor's cache, whether the memory line is shared by a number of processors, or whether the memory line is currently cached by any processor.\nA broadcast-based coherency protocol employs no tags. Instead, in a broadcast-based coherency protocol, each of the caches monitors (or snoops) requests to the system. The other caches respond by indicating whether a copy of the requested data is stored in the respective caches. Thus, correct ownership and usage of the data are determined by the collective responses to the snoops."}
{"text": "As a method of increasing a capacity of an optical disc, there are a method of increasing a recording density and a multilayering method. Further, as a method of increasing the recording density, there are a method of decreasing a channel bit length, that is, a mark length, and increasing the density in a linear-density direction and a method of narrowing a track pitch. However, if the density is increased in the linear-density direction, a problem in which inter-symbol interference is increased occurs. Further, if the track pitch is narrowed, leak-in (adjacent track crosstalk) of information from an adjacent track is increased, and quality required for a tracking error signal and a characteristic required for a tracking servo are degraded.\nMethods for reducing the adjacent track crosstalk (hereinafter simply referred to as “crosstalk”) have been proposed. Techniques of reducing influence of the crosstalk by spatially splitting reflected light from an optical recording medium into three in a track width direction, detecting each of the three into which the light is split, multiplying (weighting) the detection signal by a constant, and performing an addition operation are disclosed in Patent Literatures 1 and 2. Further, a technique capable of emphasizing and reproducing a reproduction signal of a small recording mark by further performing weighting in a beam traveling direction has been suggested in Patent Literature 1 as an idea.\nFurther, in the case of multilayering, in order to secure a sufficient recording light amount and a reproducing light amount for each layer, it is effective to keep coupling efficiency of an optical system high using only one beam in an outward path. A push-pull technique (abbreviated as a “PP technique”) has been known since the past as a method of detecting a tracking error using only one beam in an optical disc device which performs recording and/or reproducing on an optical disc having a groove structure on an optical disc recording surface, and an advanced push-pull technique (abbreviated as “APP technique”) disclosed in Patent Literature 3 and the like have been known as a method of suppressing an offset of a tracking error signal obtained by the push-pull technique which is caused by scrolling of an objective lens. In the case of a one-beam PP technique, the offset of the tracking error signal occurs at a boundary portion between an already recorded region (referred to as a “recorded region”) and a region in which recording is not performed yet (referred to as a “non-recorded region”). The offset in the recording boundary occurs when a reflectance distribution of the recording boundary is asymmetric to a track, and it becomes an amplitude grating with a period coarser than that of a track and diffracts a light beam collected on an optical disc.\nDue to the occurrence of the offset, a recording characteristic and a reproducing characteristic deteriorate, and in the least favorable case, a tracking servo is not performed. In order to suppress such an offset, a technique of performing focus control such that an offset amount of a tracking error signal is adjusted so that a level of the tracking error signal in a track jump state becomes an intermediate level between a peak level and a bottom level is disclosed in Patent Literature 4.\nFurther, techniques of detecting a tilt of a recording medium on the basis of light intensity of an interference region between 0th-order light and 1st-order diffracted light are disclosed in Patent Literatures 5 and 6."}
{"text": "White LEDs (light-emitting diodes) are devices that emits white light in combination of a semiconductor light-emitting device and a phosphor and a typical known example thereof is a combination of a blue LED and a YAG (Yttrium Aluminum Garnet) yellow phosphor. However, YAG phosphors, which contain a smaller amount of red light-emitting component, have problems that they are inferior in color rendering properties for use in illumination application, and lose their color reproducibility when used as a backlight for image display devices such as liquid-crystal display devices.\nAs a method for improving the color rendering properties of white LEDs, proposed was a method of compensating the red component by using a YAG phosphor and a red light-emitting nitride-based phosphor in combination (see Patent Document 1). It is also known that, among red light-emitting nitride-based phosphors, CaAlSiN3-based nitrides or oxynitride materials that contain an inorganic compound having a crystalline structure identical with CaAlSiN3 crystalline phase as the host crystal and additionally an optically active element, in particular Eu, as the emission center emit orange to red light particularly at high brightness (see Patent Document 2).\nIt is possible to produce the CaAlSiN3-based nitride phosphor by using a mixture of the nitrides of the elements constituting the phosphor (see Patent Document 3). In the production method for a phosphor described in Patent Document 3, a phosphor is produced by mixing nitrides of the elements constituting the phosphor and calcining the mixture obtained in a calcination furnace, using nitrogen gas as the atmospheric gas."}
{"text": "Carbon materials are light in weight and excellent in chemical and thermal stability, and moreover, they show good thermal conductivity and electrical conductivity even though they are non-metal materials. However, their uses are limited in, for example, semiconductor manufacturing processes because they have the characteristic of dust emission.\nIn view of the problem, it has been proposed that, as shown in Patent Documents 1 and 2 listed below, a chromium carbide layer composed of Cr23C6 is provided on the surface of a carbon substrate by treating the carbon substrate with a chromium halide gas. However, when a chromium carbide with a high composition ratio of chromium, such as Cr23C6, is formed, the hardness of the chromium carbide becomes high, so chipping tends to occur easily in handling or the like. As a consequence, particles are produced, so a problem with the material is that the material cannot be used for the members for which their characteristic of emitting less dust are considered important, such as the members (for example, jigs) used in a semiconductor manufacturing apparatus."}
{"text": "1. Field of the Invention\nThis invention relates to weights and a support system for exercising the abdominal muscles.\n2. Prior art\nStrengthening and toning the abdominal muscles is a common goal among exercise enthusiasts. Most of the methods used to meet this goal, however, have either biomechanical or physiological limitations.\nSuch methods that have physiological limitations are common trunk curl exercises that allow for a high number of repetitions, which have diminishing effects on the desired adaptations, i.e., stronger, firmer muscles.\nIn order for a muscle to become stronger and firmer, it has to work at certain intensity levels, below which negligible adaptations will occur. The muscles can only contract and perform for short periods of time when working at these higher intensity levels. To effectively train a muscle group for strengthening benefits, it is better to work at higher intensity levels and shorter durations.\nAnother physiological limitation of these exercises concerns neck muscle fatigue. As the various trunk curl exercises are performed, the trunk and head are lifted up and off the floor or bench. The weight of the head is supported by the neck muscles while the exercises are being performed. If the exerciser holds his/her hands behind his/her head, some weight may be relieved from the neck.\nWith many abdominal exercises, the exerciser's neck muscles become fatigued due to the prolonged and repeated isometric contractions that occur from having to support the head. This can result in stiff and/or sore neck muscles.\nThe basic abdominal curl exercise and its variations have a significant biomechanical limitation: the propensity for stretching the neck extensor muscles beyond their normal range of motion. This excessive flexing, which occurs when the exerciser allows his/her head to flex forward when performing the trunk curl exercise, can overstretch and strain the neck extensor muscles. This problem can be accentuated when the exerciser places the hands behind the top of the head and pulls, thus causing excessive neck flexion.\nAnother critical limitation of abdominal exercises concerns the effective and comfortable application of external resistance. By increasing the resistance that the abdominal muscles have to overcome to lift the trunk up and off the floor/bench when performing abdominal curl exercises, the intensity of the abdominal muscle effort is also increased. This increased effort is beneficial in that it results in a greater strengthening and toning effect on the abdominal muscles.\nThe relative degree of muscle strengthening that occurs during an athlete's training depends on the resistance which these muscles work against. Accordingly, any increase in resistance by means of weighted training devices, to impose more demands on these muscles and thereby strengthen them, is known to have a beneficial effect.\nUp until now, the methods most commonly used to add resistance to abdominal curl exercises have involved barbell plates or dumbbells placed on the exerciser's chest or behind the head. These methods are not only awkward, cumbersome, and uncomfortable, but can be unsafe."}
{"text": "This application is based upon and claims priority from Japanese Patent Application 8-077939 filed Mar. 29, 1996, the contents of which are incorporated herein by reference herein.\n1. Field of the Invention\nThe present invention generally relates to two-way communication systems which transfer information, such as search conditions and search results, between a client and a server in response to requests from the client. Namely, the present invention relates to, for example, a local-file-transfer method and a local-file-transfer system used in a WWW (world wide web) client-server system. The present invention particularly relates to a local-file-transfer method and a local-file-transfer system used in the WWW client-server system which allows large quantities of data to be transferred from a WWW client to a WWW server connected through the Internet.\n2. Description of the Related Art\nKeeping pace with recent developments of the Internet, WWW client-server systems have been increasing in the number of business application fields in which they can be used. Various execution and search requests issued from a WWW client to a WWW server involve sending large quantities of data, especially in the case where a search request composed by the client and to be sent to the server requires a significant amount of data for its definition. What is needed is a local-file-transfer method and a local-file-transfer system which allow such massive data transfer from the WWW client to the WWW server.\nSpecial types of search request information, such as genetic information, provided by a server device requires a long search-condition clause to be entered at a client device. Such a long data sequence, is generally not suitable for keyboard data entry. What\"\"s needed is the ability of storing information in a file in the client device in advance, and transferring the information to the server device at the time of the search, so that the server device can carry out the search for requested information based on the transferred file.\nTypically, data transfer from a client device to a WWW server connected through the Internet requires data entry into the client device, and such data entry is typically carried out using the following methods.\nA first method of data entry is keyboard entry of the data by an operator working at the WWW client device. The data is entered into an input field on the display which typically displays an HTML (HyperText Markup Language) menu transferred from the WWW server, and, then, the entered data is sent to the WWW server through the Internet.\nA second method of data entry is to copy (cut and paste) data to an input field from a window on the same display. In this case, the window shows a list of candidate data for transfer to the WWW server.\nThe first method of keyboard data entry into the input field is prone to key-touch errors during lengthy data entry, and is also time consuming.\nThe second method of copying (cutting and pasting) data into the input field from a window on the same display is time consuming in terms of copy-range selection, and, also, is susceptible to possible erroneous character transformation during the copying process.\nAccordingly, there is also a need in WWW client-server systems for a local-file-transfer method and a local-file-transfer system which allow massive data transfer from a client to a server without imposing an onerous work load on an operator.\nAlso, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which allow data to be transferred from a client to a server by identifying data on a client display, extracting a relevant file from a client memory device, and sending the relevant file to the server, to request the server to register the file, to execute the file, or to search registered files.\nFurther, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which allow an execution-request transfer from a client to a server to be made by entering the execution request on the client display and sending the request to the server, or by identifying the execution request on the client display, extracting file-stored execution requests from a client memory device, and sending the file to the server.\nMoreover, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which allow search-condition transfer from a client to a server to be made by entering the search conditions on the client display and sending the search conditions to the server, or by identifying search conditions on the client display, extracting file-stored search conditions from a client memory device, and sending the file to the server.\nAlso, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which can issue a completion message after completion of a data transfer from a client to a server so as to inform the server of the end of the transfer.\nFurther, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which allow a server to store client-sent information, to carry out client-issued requests based on client-registered information provided on a file-unit basis, and to allow the server to notify the client of the execution results.\nMoreover, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which allow a server to generate search parameters based on client-sent search conditions to search for client-registered information provided on a file-unit basis.\nAlso, there is a need in the WWW client-server system for a local-file-transfer method and a local-file-transfer system which can request the server to eliminate transferred information when execution results notified by the server are satisfiable.\nAccordingly, it is a general objective of the present invention to provide a local-file-transfer method and a local-file-transfer system which can satisfy the needs described above.\nIt is another and more specific objective of the present invention to provide a local-file-transfer method and a local-file-transfer system which allow data transfer from a client to a server by identifying data on a client display, extracting a relevant file from a client memory device, and sending the relevant file to the server, in order to request the server to register the file, to execute the file, or to search registered files.\nIt is a further objective of the present invention to provide a means for preparing lengthy data in advance at the client display and then transferring such data from the client to the server in a way that significantly reduces the user\"\"s work load.\nIn order to achieve the above objects, the present invention provides a method and system of transferring information from a client to a server in a client-server system which achieves significant reduction in operator work load in comparison with conventional keyboard entry. According to the present invention, the data to be transferred to the server is prestored in a file in the memory media of the client. A display is transferred to the client from the server, the display providing fields for entering data. The file in the client\"\"s memory media is subsequently entered on the display provided from the server. The file is transferred to the server at the time of search. The server device carries out the search for requested information based on the transferred file data. A search query can be constructed at the client based on combined search conditions made up from search-condition file names and search condition data previously transferred to the server.\nAnother aspect of the present invention is that a file name of a file is indicated on a display which is created based on menu information provided from the server. After indicating that file-unit-based information is to be registered; and reading the file corresponding to the file name from the memory media, the file is then transferred to the server. The above two steps are carried out when a file extracted from the memory media is transferred from the client to the server.\nAs another aspect of the present invention, an execution request is identified at the client on the display which is provided from the server, the execution request being made to the server. The execution request can be directly transferred to the server. But when the execution request is longer than a predetermined data length, the execution request is first read from the memory media and then transferred to the server.\nAs another aspect of the present invention, a search-condition identifier is entered at the client on the display provided from the server. A file is read from the memory media, and transferred to the server. This file stores the search conditions to be transferred to the server and corresponds to the search-condition identifier that was indicated.\nAs a further aspect of the present invention, a file-transfer-completion message is issued from the client when file transfer to the server is finished.\nAs a further aspect of the present invention, the server allocates file space for storing the file transferred from the client. The file-unit-based information transferred from the client is stored in the server\"\"s file space. A receipt-completion message to the client is issued from the server when all files transferred from the client are received. The server carries out the execution request or search conditions request obtained from the client. Upon completion, the server then sends the execution results or search results to the client.\nAs a further aspect of the present invention search conditions are entered on the display of the client. This display was previously created based on menu information provided from the server. The search conditions to be used by the server are transferred from the client to the server. Search parameters are generated at the server based on the search conditions obtained from the client. Searching is accomplished at the server based on the generated search parameters for file-unit-based information transferred from the client in advance. The search results are transferred to the client.\nAs a further aspect of the present invention, an elimination request is sent from the client to the server based on one of the execution results and search results which had been previously sent from the server to the client. A file transferred from the client is then expunged at the server based on the elimination request sent from the client.\nIn order to achieve the above objectives, the present invention provides a system for transferring information between a client and a server which includes: a local-memory for pre-storing information, the information having a file structure for transfer to the server and a search-condition transfer unit for reading at least one file from the local-memory unit and for transferring this file to the server. The server has a search unit which executes a search based on search conditions contained in this file transferred from the client and transfers the results back to the client.\nAs a further aspect of the present invention, the search unit of the server has a search-condition storage unit for storing at least one file transferred from the client; and a search-execution unit for transferring search results to the client. The search results are obtained from a search based on combined search conditions made up from at least some of the files stored in the search-condition storage unit.\nAs a further aspect of the present invention, the search-execution unit includes a search-result transferring unit for transferring search results to the client. The search results are obtained from a search based on combined search conditions made up from at least some of the files and search-condition data transferred from the client.\nAs a further aspect of the present invention, the system further includes a file-transfer instructing unit, a file indicating unit and a file-transfer unit. The file-transfer instructing unit accepts an entry of data indicating that file-unit-based information is to be transferred. The entry of the data is made on a display which is created based on menu information provided from the server. The file indicating unit accepts at least one file name. The file-transfer unit reads at least one file indicated by the file indicating unit from the local-memory unit and transfers this file to the server.\nAs a further aspect of the present invention, the client of the system further includes an execution-request instructing unit which accepts an entry of an execution request on the menu information provided from the server. The client also further includes an execution-request-transfer unit which transfers the execution request indicated by the execution-request instructing unit to the server. Alternately, the execution-request-transfer unit reads the execution request indicated by the execution-request instructing unit directly from the local-memory and then transfers the execution request to the server.\nAs a further aspect of the present invention, the client of the system further includes a search-condition indicating unit which accepts an entry of a search-condition identifier on a display created based on menu information provided from the server. A search-condition-file-transfer unit of the client reads the files from the local-memory unit based on the search-condition identifier indicated by the search-condition indicating unit and transfers the files to the server. A search-condition-transfer unit transfers the search-condition identifier indicated by the search-condition indicating unit to the server.\nAs a further aspect of the present invention, the client of the system further includes a transfer-completion-message issuing unit for sending a file-transfer-completion message to the server when either the file-transfer unit, the execution-request transfer unit, or the search-condition-file-transfer unit finishes transfer of all files to the server.\nAs a further aspect of the present invention, the server of the system further includes: an execution-request-storage unit for storing files of execution requests sent from the client and an execution unit for transferring execution results to the client. The execution results are obtained from an execution based on combined execution conditions made up from a plurality of the files of the execution requests stored in the execution-request-storage unit.\nAs a further aspect of the present invention, the execution unit of the system further includes an execution-results-transfer unit for transferring execution results to the client, the execution results being obtained from an execution based on combined execution conditions made up from the files of the execution requests and execution-request data transferred from the client.\nAs a further aspect of the present invention, the server of the system further includes a completion-message issuing unit for issuing a receipt-completion message when the server receives all files transferred from the client.\nAs a further aspect of the present invention, the server of the system further includes a parameter generating unit for generating parameters based on the search conditions obtained from the client. The parameters are used for searching in an information-storage unit of the server. The server includes a search unit for searching in the information-storage unit based on the search parameters generated by the parameter generating unit.\nAs a further aspect of the present invention, the client of the system further includes an elimination-request issuing unit for sending an elimination request to the server based on one of the execution results and search results which are sent from the server to the client; and an elimination unit for eliminating the files stored in the information-storage unit of the server based on the elimination request sent from the client.\nAccording to the present invention, transfer of lengthy data from the client to the server does not require the operator to make keyboard entry of the data. File-unit-based information corresponding to either data registration, an execution request, or a search request is extracted from the memory unit of the client to be sent to the server. Thus, only key information needs to be entered in order to read desired data from the memory unit. Also, the server device stores the search request sent from the client in the memory unit so as to carry out the search request made by the client at a later stage.\nKey information for reading file-unit-based information to be transferred to the server from the client memory unit is entered in a menu display provided by the server using predetermined protocols. Since the file-unit-based information is transferred to the server, keyboard entry of lengthy data is unnecessary.\nThe execution request itself is entered at the client and subsequently transferred to the server. Alternately, the long-data-length request stored in the local memory unit is read according to the execution request and transferred to the server. In the case of a short data length which imposes no load on the operator, the execution request itself entered on the display can be transferred to the server. Alternately, with a long data length imposing the load on the operator if working manually on the data entry, the file-unit-based information extracted from the local memory unit according to the input execution request (local-memory-means-search key) can be transferred to the server. Namely, a request having a short data length can be transferred to the server without resorting to the extraction of the file-unit-based information from the local memory unit, whereas a request having a long data length can be extracted from the local memory unit to be transferred to the server.\nSearch conditions entered on the display provided by the server are transferred to the server as long as the search request has a short data length. Alternately, long search conditions stored in the local memory unit are extracted according to the input search conditions and transferred to the server if the search conditions have a long data length. Either one of these two operations can be selected for transfer to the server. Accordingly, simple search conditions are directly transferred to the server, and long and complex search conditions are extracted from the local memory unit to be transferred to the server since the keyboard data entry would impose a heavy load on the operator.\nThe client sends a file-transfer message to the server when the data transfer from the local memory unit of the client to the server is completed. In response to the file-transfer message, the server executes various requests indicated by the client.\nThe server stores the files sent from the client, and sends a receipt-completion message to the client after receiving all the files from the client. Then, the server carries out a request from the client, and sends execution results to the client. Based on the execution results from the server, the client can make a decision as to whether to issue a further request, to terminate a request process, or some other operation.\nThe server generates parameters indicating storage positions of files sent from the client, so that the server can access to the starting positions of the stored files.\nWhen the execution results are sent from the server to the client, the client makes a decision based on the execution results, and sends an elimination request for purging files when there is no further request or when files stored in the server are no longer necessary. In response, the server can eliminate the files indicated by the elimination request, so that the server can avoid a waste of the file storage.\nThe client reads a file from the local memory unit, and sends it to the server. The server carries out a search based on the search conditions contained in the file, and sends the search results to the client.\nSearch is conducted based on combined search conditions made up from the search conditions of a plurality of files stored in the search-condition storage means.\nThe server carries out a search based on combined search conditions made up from both the search conditions of a plurality of stored files and the search-condition data sent from the client.\nThe client indicates a file transfer operation and a relevant file name so as to transfer a file to the server.\nThe client transfers to the server an execution request which is entered manually. Alternately, the client reads an indicated execution request from the local memory unit to transfer it to the server. The client transfers to the server search conditions which are entered manually. Alternately, the client reads indicated search conditions from the local memory unit to transfer them to the server.\nThe client sends a file-transfer-completion message to the server when the file-transfer unit, the execution-request-transfer unit, or the search-condition-file-transfer unit finishes transfer of all the files. In response, the server can start executing a process such as a search process.\nThe server can carry out execution of a plurality of execution files sent from the client by combining them together.\nThe server can carry out an execution based on a combined execution request made up from the execution-request data sent from the client and the execution-request files.\nThe server sends a receipt-completion message to the client when the server has finished receiving all the files transferred from the client.\nThe server searches for stored files based on parameters generated for the search.\nThe client makes a decision based on the execution results sent from the server, and sends an elimination request for purging files when there is no further request or when files stored in the server are no longer necessary. In response, the server can eliminate the files indicated by the elimination request, so that the server can avoid a waste of the file storage. A summary of the improved results of the invention are as follows.\nA large amount of data transfer from the client to the server is achieved in the WWW client-server system. Search for files sent from the client is conducted based on combined search conditions made up from a plurality of search-condition files.\nSearch for files sent from the client is conducted based on combined search conditions made up from both the plurality of search-condition files and the search-condition data. Information to be transferred to the server is indicated on a display, and a file extracted from the client memory means according to the indication is transferred to the server.\nThe execution request such as programs to be executed by the server is entered on the client display, and is transferred to the server. Alternately, the indicated execution request is extracted from the client memory unit, and is transferred to the server.\nSearch conditions to be used by the server are entered on the display, and are transferred to the server. Alternately, the indicated search conditions are extracted from the client memory unit, and are transferred to the server.\nThe client sends a completion message to the server when all the transfer from the client to the server is completed, so that the server can know the end of the transfer.\nThe server can carry out an execution based on combined execution conditions which are made up from a plurality of stored execution files such as those containing programs.\nThe server can carry out an execution based on combined execution conditions which are made up from both the plurality of stored execution files and the execution-request data.\nThe server stores information sent from the client and carries out a request from the client based on the stored information and then sends execution results to the client. The server receives search conditions from the client, and generates parameters to be used for a search of the client-sent information based on the search conditions.\nWhen the execution results sent from the server to the client are satisfactory or when files transferred to the server are no longer necessary, the client can make a request to the server to purge the transferred files.\nAs described above, according to the present invention, the local files indicated at the client device are sent to the server. This way, data registration to the server database through a network such as the Internet or a data-base search at the server device can be readily achieved although registered data or search keywords are lengthy.\nEspecially, when a search is made through a network for gene arrangements stored in a gene-arrangement data bank, the present invention can achieve significant reductions in the operator\"\"s work load and work time compared to conventional keyboard entry of gene-arrangement data up to several kilo bytes.\nOther objects and further features of the present invention will be apparent from the following detailed description when read in conjunction with the accompanying drawings."}
{"text": "As a result of progress in current CMOS (complementary metal oxide semiconductor) technology, such as advances in 65-nm and 45-nm processes, in a programmable logic device (PLD), such as a field programmable gate arrays (FPGAs), have increased logic elements and I/O buffer (or pin) densities. Because simultaneous switching noise (SSN) is associated with the number of toggling I/O buffers, signal voltage level, and the switching rate of I/O buffers, it is important for FPGA users to be able to evaluate SSN performance in a “chip-package-board” environment. The SSN caused by FPGA output buffers is widely known as simultaneous switching output (SSO) noise and can be differentiated from the SSN caused by input buffers.\nSSN in an FPGA system may be attributed to two primary factors: the mutual inductive coupling among switching I/O buffers and the impedance profile of a power distribution network (PDN) including die, package, and printed circuit board (PCB). In essence, reducing SSN is a design-cost issue. While designers can minimize the mutual inductive coupling by increasing the ratio of FPGA ground pins (or return-current pins) to I/O buffers, this approach sacrifices I/O buffer densities. Engineers can improve PDN performance by increasing on-die capacitance and adding on-package decoupling capacitors, but this approach increases costs.\nBecause FPGAs are programmable, they fit into a wide variety of user applications, so it is useful for designers to have a tool to determine their own SSN budget without additional costs. Ideally, this kind of tool requires instantaneous and accurate result predictions for various FPGA I/O buffer assignments and ideally would provide an optimum buffer assignment under certain design constraints. In the past, designers have constructed system-level, “SPICE-like” models to anticipate SSN in FPGA systems. These models are based on an understanding of SSN cause mechanisms and correlate well with bench measurements, helping IC and packaging designers improve designs.\nHowever, these models are so complicated that they require signal/power integrity expertise to perform time-consuming, system-level simulations. It is cumbersome for all FPGA designers to perform the same level SSN analysis without considering their different design margins. Therefore, there is a need for a tool to help FPGA users execute a comprehensive FPGA SSN analysis in a short design cycle, though the invention claimed below has applicability to other applications beyond this particular application, as will become apparent from the following description and the drawings which accompany it."}
{"text": "This invention relates to a water treatment unit for the circuit water of open-circuit cooling systems, more particularly for cooling tower circuits, comprising a pipe with a filter unit connected to the circuit water stream of the cooling circuit for branching off and returning a part stream of the cooling water and a pipe for adding fresh water connected to the circuit water stream.\nWater treatment units are known, for example from EP 0305897 (Henkel KGaA). The circuit water is passed through heat exchangers where it is heated and then pumped to the cooling tower. Here, the water is cooled by evaporation in an air stream and is re-used for cooling plants, reactors, etc. With open cooling systems such as these, the water lost through evaporation has to be replaced. In addition, additives for preventing corrosion of metal components and for preventing the formation of organic, inorganic or biological deposits have to be added to the cooling circuit water. It is known that the evaporated water can be replaced by untreated water, for example drinking water or ground water, or by pretreated water, for example decarbonized or deionized water. Through the addition of salt-containing water, i.e. water containing calcium or magnesium chloride, the salts are concentrated in the cooling system and generally lead to deposits of calcium and magnesium, particularly at places where the temperature is relatively high. In addition, relatively high chloride levels promote corrosion. It is known that part of the cooling water can be periodically xe2x80x9cflooded offxe2x80x9d to keep the salt concentration within tolerable limits.\nIn order to reduce the concentration of suspended solids, it is known that part of the cooling water, generally 5% of the circuit water, can be removed from the circuit as a part stream and passed through filters where the suspended solids are removed from the cooling circuit water. The filters are periodically cleaned by backwashing with the cooling circuit water, the backwash water being discharged as wastewater. The backwash water contains only about 1% of solids. As a result of the above-mentioned flooding off of the cooling water and the discharge of the backwash water as wastewater, water treatment chemicals are lost and have to be replaced. With relatively large cooling systems, this requires an elaborate measuring and control system as described, for example, in the above-cited EP 0305897 B1.\nIn addition, in view of legal requirements and the sensitive chemical environment in the cooling circuit, particularly so far as corrosion and deposits are concerned, the cooling circuit water has to be frequently monitored by chemical analysis. Known systems operated with low-salt or salt-free water are not free from wastewater either because the filters are cleaned by backwashing and the backwash water is discharged. As a result of this loss of filter backwash water, therefore, the cooling circuit water again has to be frequently analyzed and water treatment chemicals continuously added.\nAccordingly, the problem addressed by the present invention was to provide a unit of the type mentioned at the beginning for the treatment of circuit water of open circuit cooling systems and a corresponding treatment process which would satisfy the following requirements. Suspended solids would have to be largely removed from the circuit water to prevent the deposits and reduce corrosion. This applies in particular to solids with an abrasive effect. The unit would have to lend itself to operation with low-salt and, more particularly, salt-free water so that few, if any, corrosion-promoting substances, for example chloride ions or sulfate ions, would be present in the cooling circuit water. Wastewater-free operation would have to be possible so that water treatment chemicals would not be discharged. On the one hand, this would reduce operating costs and, on the other hand, would avoid the introduction of chemicals into receiving waters. The quantity of water to be added would have to be minimal; in particular, only the water evaporated in the cooling tower would need to be replaced. This requirement would be fulfilled in particular by wastewater-free operation of the unit. Particularly stable operation of the cooling system would have to be possible in order significantly to reduce the expense of monitoring by chemical analysis in continuously and discontinuously operated systems. Finally, the need for forced circulation during stoppages as necessary in the prior art would have to be eliminated."}
{"text": "1. Field of the Invention\nThe present invention relates to a film projection system and more particularly to a film projection system wherein in response to the actuation of a selector system, a cartridge containing a film program is automatically transferred from a storage magazine to a film reproduction unit in which the film program is reproduced, the cartridge being returned to the storage magazine upon completion of the film program.\n2. The Prior Art\nThere are many types of film projector systems in the prior art. In one such system, a storage reel containing a film program must be manually placed on the projector and the film manually positioned or threaded through the projector to permit reproduction of the program retained on the film. More recently, projector systems have been developed in which, after the storage reel is placed in the projector, an end of the film is inserted into an automatic threading device and the projector will be threaded automatically. In all of these prior art projection systems the storage reel containing the film has to be positioned on the projector by hand. In addition, the film must either be threaded through the projector by hand or one end of the film must be inserted by hand into an automatic threading device in the projector. In either case, the presence of an individual to assist in the operation of the projector is required which severely limits the use of the projection system, especially in the field of entertainment.\nSound systems have been developed such as the wellknown record player or \"Jukebox\" which permit individuals to select a sound recording from a collection of recordings, and the sound recording then played automatically. While such sound systems have enjoyed tremendous popularity, no such system has been developed applying similar principles to reproduce film programs in which the individual would not only hear the sound recording but would also see a visual presentation of the artist performing. Such a system would open vast areas to pre-programed film performances which have heretofore been out of reach because of the inability to select and play automatically, a pre-recorded film. Presently, individuals viewing film programs in association with sound recordings must rely on the person in control of or who operates the projection system to select or choose the film to be reproduced.\nIt is toward the elimination of these and other problems that the present invention is directed."}
{"text": "Lately with the development of the housing industry, the secondary housing industry has progressed very rapidly, including manufacture of built-in furniture, kitchen sinks, washing stands as unit equipment.\nKitchen sinks and washing stands are one of the interior fixtures essential for every day life, and a variety of such products directed to the taste of customers are being presented on the market.\nSince kitchen sinks and washing stands are fixtures which inevitably involve use of water, a top which supports a basin is provided horizontally as one of the constituent members. As such a top there is usually employed that having a structure obtained by making the surface of a wood derived core material repellent to water and at the same time giving a color so as to serve also as surface decoration or bonding a melamine decorative board having a woodgrain pattern. The kitchen sinks and washing stands are, as described above, fixtures utilizing water indoors, and therefore the used water should not leak into other areas than the sinks and washing stands.\nAccordingly, various devices have been imparted to the aforementioned tops fitted with basins in order to prevent leakage of water spilt on the surface of the top through a joint part between a flange part of the basin and the edge of a hole in the top for imbedding the basin. One such example is a product in which a basin flange part is mounted on a top, the joint part thereof with the edge of a hole in the top for imbedding the basin is sealed by fixing a rubber packing or other sealing materials, and the top and the basin flange part are fixed together by screwing from the rear surface of the top. However, those having such a structure suffer from disadvantages that with prolonged use, the rubber packing or other sealing materials are degraded and allow spaces around the edge of the aforementioned imbedding hole to cause water leakage, thereby resulting in swelling and corrosion of the top itself, which then forms gaps between the surface of the top and the basin flange part, finally permitting water spilt on the top to drop onto the floor instead of flowing into the basin, thus making the neighborhood of the sink or washing stand very moist.\nAs an alternative to those having such structure that the joint part between the basin flange part and the edge of the hole in the top for imbedding a basin is sealed using a rubber packing or other sealing materials, there is a structure in which a hole in the top is made larger than the basin itself, and the top decorative board alone is provided with a hole having the same dimention as the open part of the basin, and the basin flange part is fitted and directly adhered to the hole edge of the decorative board on its rear surface using a synthetic resin. By being of such a structure, since the flange part of the basin is situated lower than the surface of the top, all the water spilt on the surface of the top flows into the basin, thus solving the aforementioned problem of drainage of water spilt on the surface of the top. However, with this structure, the strength of the bonded part between the top decorative board and the basin flange part is weak, and it therefore tends to be damaged during transport and the decorative board is also easily impaired. Moreover, the finish treatment of the synthetic resin used for adhering the joint part between the hole edge of the decorative board and the basin flange part is troublesome and requires a lot of labor and time, and therefore it is not applicable to manufacture of products having a structure suitable for mass production.\nAs mentioned above, in fabricating tops fitted with basins used in kitchen sinks, washing stands etc., various problems are present from a standpoint of working to prevent water leakage in the joint part between the basin flange part and the hole in the top for imbedding a basin.\nThe present invention aims to present tops fitted with basins which can provide good water drainage and have a structure permitting their mass production by eliminating problems involved in the structure and fabrication of tops fitted with basins which are constituent members for kitchen sinks and washing stands, as well as a process for their production."}
{"text": "The term “displacement element” is to be understood comprehensively in the present case. It includes hatchbacks, trunk lids, engine hoods, doors, in particular side doors, cargo space floors, or the like of a motor vehicle.\nHowever, the drive configuration under discussion is primarily used in hatchbacks and side doors in motor vehicles. It is used for the motorized displacement of the respective displacement element in the closing direction and in the opening direction. It is typically important that the drive configuration also permits manual displacement operation in addition to the motorized displacement operation. The manual displacement operation is significant in particular in case of emergency, for example, in the event of a crash or power failure.\nThe known drive configuration (DE 20 2005 007 155 U1), from which the invention proceeds, is associated with a hatchback. The drive configuration is equipped with two spindle drives, which each have, in a compact module, a drive motor, an intermediate gearing having clutch, and a spindle-spindle nut gearing. A spring configuration is provided in the respective module, which counteracts the weight force of the associated hatchback. The known drive configuration further has a drive controller, which is used to activate the two drives, in particular the two drive motors. The drives are not designed as self-locking, so that manual operation is implemented in a simple way.\nHatchbacks of significant size and/or significant weight can be adjusted by motor using the known drive configuration. This opens up new degrees of freedom in the design of such hatchbacks. However, an increased risk it event of a failure of drive components is also connected to the increase of the weight."}
{"text": "A well-recognized goal within the field of solid-state color image sensors is the production of a solid-state color imager which is highly sensitive to light, and which produces a clear image while being inexpensive to manufacture. In pursuit of this goal, a number of different types of solid-state color imagers have been produced.\nIn one such imager, panchromatic image sensing elements in an array are selectively sensitized to color by means of an integral array of color filters disposed over the array of image sensing elements. Highly efficient configurations for such filter arrays which maximize the amount of usable inforhave been described for example in U.S. Pat. No. 3,971,065 issued July 20, 1976 to Bayer and U.S. Pat. No. 4,047,203 issued Sept. 6, 1977 to Dillon. However, the resolution capabilities inherent within such arrays are limited by the number of sensing elements that can be placed on the array and are further limited in that only a portion of each element in the array contributes to the resolution of fine detail. Consequently, the spatial resolution of such integral-filter color image sensing arrays, while optimized for the particular design, will not be as high as monochrome image sensing arrays of the same number of elements.\nAnother structure as proposed within British Pat. No. 2,029,642 and Japanese Patent Applications Nos. 55-39404; 55-277772; 55-277773; and 51-95720 is constructed such that the photosensor is superimposed on top of the information transfer device or solid-state base which is capable of a switching function. The base may be an MOS or CCD switching device. Such devices are described in detail within British Pat. No. 2,029,642 the disclosure of which is incorporated herein by reference. Such structures have potentially high sensitivity due to a larger sensing area than is present within conventional imaging devices where the photosensor is located on the same level as the information transfer device. However, such devices must utilize multi-color filters and the loss of resolution is comparable to normal solid-state imagers as discussed above. In addition, to produce such a structure, the color filters must be arranged in a particular pattern on the image sensing element which creates difficulty in the alignment and bonding of the filters making the production of such devices complex and expensive.\nImprovements over the structure disclosed in the aforementioned British Patent are described in European Pat. No. 0046396 and Japanese Patent Application No. 56-133880. The improved structure in the European Patent uses three MOS elements on the solid-state base for each image sensing element. The photo-carrier generated in a photoconductive layer is applied to the gate of an amplifying MOS as compared to the drain or source in the British Patent which in turn is connected to a switching MOS transistor. The third MOS transistor provides a means for resetting the photoelectric conversion elements. Solid-state color imagers with this structure improve resolution and signal-to-noise ratios by preventing picture cells which are not being scanned from providing false signals.\nThe improved structure in the Japanese Patent Application No. 56-133880 uses a Metal Nitride Oxide Silicon Semi-conductor (MNOS) field effect transistor for detecting and storing data for photosignals, as electrical signals. Similarly to the European Patent described above, the photosignal is applied to the gate of the transistor. Solid-state color imagers with this structure give rise to a nonvolatile memory effect which results from the use of the MNOS semiconductors.\nBoth of the devices described in the European Patent and Japanese Application No. 56-133880 must still use multi-color filters arranged in a particular pattern on the image sensing element.\nA technique for eliminating color filters in a vidicon is taught in U.S. Pat. No. 3,617,753 of Kato et al. The vidicon includes a conventional semiconductor layer having a substrate on a plurality of p-n diodes which store electrical signals representing light intensity. An electron beam scans the p-n diodes to provide video read out. By stepping the thickness of the semi-conductor substrate through which the light passes to the p-n diodes, different wavelength light impinges on the p-n diodes, depending on the size of the step. In this manner different groups of p-n diodes can store different color light. Alternatively the p-n diodes can be formed at varying depths from the surface, thereby effectively stepping the thickness of the substrate. In another embodiment, solid state scanning can be provided instead of electron beam scanning. There a junction device and a MOS element is provided at each pixel and selective etching of the substrate results in varying distances between the light receiving surface of the semiconductor substrate and the junction device of the pixel. The apparatus disclosed is not planar due to the stepped or cutout arrangement and does not have the advantage provided by systems using photoconductors as the light responsive element.\nA solid-state color image sensing array has been developed wherein the potential resolution is equal to that of a monochrome array of the same size. Such a sensing array has a plurality of superimposed channels (e.g., three superimposed channels for a three-color device) wherein each channel has a different spectral response due to differential absorption of light by a semiconductor material. (See Research Disclosure, August 1978, Vol. 172, Disclosure No. 17240 entitled: \"Color Responsive CCD Imager\" available from Industrial Opportunities, Ltd., Honeywell, Havant, Hampshire P091EF, U.K.) However, extremely complex and expensive processes are necessary to produce such devices due to the necessity of superimposing the three channels. When utilizing the CCD (charge-coupled device), the channels which carry the information signal must be carefully constructed within precise limitations making the construction complicated and expensive. Although it is possible to produce a single channel on a substrate, it is complicated and difficult to superimpose additional channels thereon.\nDevices such as those described in Disclosure No. 17240 indicate that it is possible to produce multiple superimposed varied channels in silicon crystal which can act as multi-channel superimposed color-sensing devices. However, in addition to the expense and complication of their manufacture, as mentioned above, the color separation and selectivity of these devices is poor due to the inherent limitations of the materials used. The materials used in making such devices act as CCD channels which must have good single crystalline properties as well as color selective photosensors.\nAs mentioned above, there exists a need within the field for a solid-state color imager which is highly sensitive to light and which gives sharp, detailed resolution of the image. By utilizing a device wherein the multi-colored filters are superimposed over the image-sensing elements in an array, the resulting image, as described within U.S. Pat. No. 3,971,065, has limited resolution capabilities, limited sensitivity and is complicated and expensive to produce due to the necessity of precisely placing the multi-color filters. Increased sensitivity can be obtained by utilizing a device wherein a photosensor is superimposed on top of the information transfer device, as described within British Pat. No. 2,029,642, European Pat. No. 0046396 or Japanese Application No. 56-133880. However, resolution of such devices is still somewhat limited because they require the use of multi-color integral filters which also increases the complexity and expense of their production. By utilizing a device having a sensing array of a plurality of superimposed channels, it is possible to obtain a resolution equal to that of a monochrome array. However, complex, expensive manufacturing techniques must be utilized to superimpose three channels on top of each other.\nThe present invention utilizes a plurality of photosensitive layers which are superimposed on each other and over the base which includes image sensing elements which are used for detecting a single color of light. The invention eliminates the need for multi-color integral filters since each photosensitive layer and the elements on the base of each detects a different color of light. The device has a resolution equal to that of a monochrome array of the same size and can be produced by simple, conventional, inexpensive techniques."}
{"text": "The present invention relates to a carburetor chamber for a power-driven machine, particularly a chain saw, which is provided, at its top, with an engine cylinder and, at its lower part, with an engine crank case and, at the rear part of the cylinder, with a carburetor in a manner that the carburetor is juxtaposed with the cylinder.\nGenerally, in a machine arranged to operate with the use of an internal combustion engine as a power source, a connecting or communicating of an air-fuel mixture outlet of a carburetor with a suction port of a engine cylinder is effected through a bellows type gas pipe in consideration of the thermal expansion, etc.\nIn this case, however, the fuel in the air-fuel mixture is very likely to attach to and stay on an inner surface of the bellows, so that the air-fuel ratio is varied due to a dropping of such fuel, so that unexpected variations in revolution occur during an idling speed of the engine. This becomes a cause of the troubles.\nFurther, in the chain saw, an elongate flexible pipe and the like are used to feed chain-lubricating oil under pressure. For this reason, for example, the problems of pipe breakage and oil leakage are liable to arise as the particular problems of the chain saw. Further the inconveniences such as that of making it difficult to adjust or control the engine output and the engine idling speed due, for example, to a displacement in attachment position of the throttle valve operation lever of the carburetor, have hitherto arisen."}
{"text": "1. Technical Field\nThe present invention relates in general to a method and system for permitting rapid implementation of new services within communications systems. In particular, the present invention relates to a method and system for taking direct advantage of new technological advances in network design and traffic transport. More particularly, the present invention relates to a method and system for allowing independence of session processing from the RF technology utilized in the wireless or satellite system, thus maximizing the re-use of session processing functionality, software, and hardware in diverse wireless and satellite networks. Still more particularly, the present invention relates to a method and system for partitioning the functions of RF channel assignment, session admission control, authentication, subscriber profile management, fast message services, billing operations, and OAM (Operations, Administration and Maintenance), thus simplifying the design of each function and maximizing the re-use of existing components, thereby allowing network operators to control session admission, session setup and release, the type and number of sessions that concurrently exist, and the quality of service for the sessions.\n2. Description of the Related Art\nFactors driving the need for wireless and satellite communications arise from changing user needs and demands. Previously, public network needs were driven primarily by voice data telephony. Wireless and satellite communications are sometimes referred to as multiple access communications systems, reflecting the fact that they involve multiple communications terminals that share the same communication medium. Multiple access communications systems are distinguished from so-called \"point-to-point\" communications which are found in terrestrial telephone networks in which physically separate communications links are utilized.\nData traffic has grown slowly until recently. With the lower cost in telecommunications and increase in processing power of computers, the number of users accessing wireless and satellite communications networks has increased. The needs of these users include, for example, video telephone, low cost video conferencing, imaging, high definition television (HDTV), and other applications requiring multimedia data transfers. Multimedia combines different forms of media in the communication of information between a user and a data processing system, such as a personal computer. A multimedia application is an application that utilizes different forms of communications within a single application. Multimedia applications may, for example, communicate data to a user on a computer via audio, text, and video simultaneously. Such multimedia applications are usually bit intensive, real time, and very demanding on communications networks.\nA multiple access communications system is comprised of a set of transmission and reception devices that include individual communications terminals, such as a personal computer or gateways which act as communications hubs for multiple terminals. A communications session is initiated when a communications request is delivered from an originating terminal to some form of Network Control Center (NCC). In a satellite communications system, a connection request is first delivered to a transceiver aboard a satellite. Thereafter, the satellite then relays the connection request message to the NCC which is comprised of the various control devices and entities that must process the request before the connection will be permitted to proceed.\nBefore the NCC will authorize the satellite to perform the requested connection, the control devices and entities, referred to hereinafter as control interfaces must perform various functions such as user authentication, billing updates, RF channel allocation and updates, etc. As previously discussed, multimedia communications within a multiple access communications system is enormously complex and the control functions required to setup and monitor such communications is likewise multifaceted and complex.\nAs an example, an NCC may include several distinct devices responsible for different control functions. For simplicity of explanation, these devices will be referred to hereinafter as \"control interfaces\". Currently an NCC utilized in a multiple access communications system may employ many different control interfaces that individually perform particular control functions such as those involved in communication setup or in monitoring a parameter, such as channel quality, after the communication link has been established. These control functions are often handled as discrete functions and are therefore performed entirely separately without adequate interfacing to permit efficient, orderly and reliable control administration.\nThe sheer complexity and the substantial potential for confusion resulting in errors are an inherent problem involved in controlling such communications sessions. One example of such an error occurs when an authentic but financially defaulted telecommunications account is erroneously permitted to place a long-distance call. Many other such examples are possible.\nBased on the foregoing, it can be appreciated that a need exists for a method and system for issuing the necessary instructions to and from communications control interfaces. Such a method and system may be utilized for coordinating activities among the various network devices in wireless and satellite communications networks. Among the advantages of implementing such a method and system is that such a technique would prevent a communications control center (located on a satellite, for example) from receiving conflicting control messages. Such a method and system would thus increase the efficiency and reliability of the control phase of wireless and satellite communications systems. Such a method and system, if implemented, would also be useful to control and regulate network traffic congestion while maintaining Quality of Service (QoS) guarantees for users of such wireless and satellite communications networks."}
{"text": "Various data measuring and collecting devices and methods are useful for analyzing a club, racket, bat, or steering wheel (herein generically referred to as “sporting devices or equipment” or “sports equipment”) during a swing or other operation. In a similar manner, the effectiveness of an impact of a ball with sporting equipment during a swing can be measured in terms of initial ball launch conditions. These launch conditions are determined principally by the velocity of equipment at impact and the loft and angle of the ball contacting surface relative to the intended trajectory of the ball's flight. Ultimately the swing of the user and the force applied on the grip by the user, determine the launch conditions of a ball. There are two general methods for analyzing the equipment during a swing; visual analysis and quantitative variable analysis.\nThe method of analyzing a swing using visual analysis typically is conducted by an instructor capable of visually discerning swing variables and suggesting corrections to the swing to provide improvement. However, not every user has ready access to professional instruction or can translate an instructor's feedback into a more efficient swing. An instructor can also not “see” quantitative factors such as force and acceleration.\nQuantitative variable analysis employs sensors to directly measure various mechanical or physical properties of the equipment during the swing motion. Sensors, such as force sensors or inertial sensors, typically are attached to the handle or the striking surface of the equipment or can be attached to the hands of the user of the equipment. Data collected from these sensors then may be transferred to a signal processor via wires or radio waves, and can be presented in various graphical formats, including graphical and tabular charts. A drawback associated with the use of existing instrumented golf clubs and other sports equipment is that the sensors and associated wires can be obtrusive to the user when the user attempts to swing the club or racket. The force and acceleration profile obtained is not then representative of the user's profile when using an unencumbered device.\nSwing characteristics will also be different between practice conditions, where a player may be relaxed and more thoughtful about a shot, and actual play where other tensions come into play.\nAn objective of the present invention is to provide an instrumented grip for sports equipment that delivers an enhanced comfort level to the user when grasped in the course of actually playing the sport or in a practice environment. The device therefore provides a means for comparing swing parameters in the idealized setting of a practice, with actual performance while playing the sport.\nA further objective is to provide a monitoring, diagnostic, and training device integrated with the grip for sports equipment. The device can be used without any interference with the natural feel and comfort that the equipment would have if the device were not there. In particular it is not necessary for a user to wear special equipment such as instrumented gloves to produce meaningful data from use of the equipment.\nA still further objective is to provide a grip that comprises a lightweight open structure for weight and weight distribution control. The open structure of the grip also allows for placement of electronic equipment in the grip."}
{"text": "Field of the Invention\nThis invention pertains in general to the field of stimulation devices, and especially to conductive implantable stimulation devices.\nDescription of the Prior Art\nTranscranial direct current stimulation (tDCS) is a non-invasive technique that can be used to modulate brain activity and to promote neuroplastic changes for therapeutic aims. Recent advances include therapeutic success in major depression treatment, potential in chronic pain treatment and stroke rehabilitation. Specific targeting of electrical current is hampered by the skull, which has poor conductivity. This is a well-known problem in brain bio-electromagnetism. The skull thickness and conductivity properties are unknown, which makes the prediction of current flow into the brain a very difficult task that require detailed measurements and modelling of electrical properties of the anatomy and structure of brain and skull.\nMoreover, in animal studies with tDCS has been shown that nearly 50% of the electrical current of tDCS does not reach the brain, but rather flows from anode to cathode on the skull 11 surface. Therefore, the output current from the device needs to be increased to reach a desired target levels in the target brain tissue. This, in turn, creates complications such as adverse sensations, skin irritation and possibly skin burns.\nIn chronic diseases, such as Parkinson's disease or chronic pain, one potential treatment approach is to implant a deep brain stimulation (DBS) device (Parkinson's) or an epidural motor cortex stimulator (MCS, chronic pain). In DBS, the electrodes are installed close to deep brain structures (subthalamic nucleus or globus pallidus) to specifically target these structures. In MCS, the aim is to reach specific cortical areas to modulate cortical activity associated with the pain relief. Installation of a DBS electrode is expensive and can lead to serious adverse effects, such as stroke, cerebral infections. Moreover, technically complicated equipment requires periodical maintenance, battery charging or changing that needs to be completed in an operation. Similar complications are present in the MCS implantation.\nThus, a need for an improved stimulation device solving or at least mitigating the above problems is needed."}
{"text": "In the prior art, as a method and device that can integrally mold a pipe with a synthetic resin, there has been known the method and device that uses a mold in which a cavity has on its one end a floating core, having a diameter corresponding to the inner diameter of the pipe, and a pressure port through which a pressurized fluid is pressure-injected, and the cavity has on its other end an outlet. After the inside of the cavity is filled with a molten resin, the pressurized fluid is pressure-injected through the pressure port to move the floating core to the outlet side, and, thus, to form a hollow in the resin in a main cavity. At the same time, an excess resin is extruded from the outlet, whereby a hollow pipe is integrally molded (Patent Documents 1 and 2).\nHowever, in the method and device described in the Patent Documents 1 and 2, the hollow of the pipe is formed as sort of a moving track of the floating core. Since the floating core moves in only one direction, there is a problem that a pipe with a branch cannot be produced. It is considered to provide a plurality of floating cores and move the floating cores in different directions so that the pipe with a branch can be integrally molded; however, because of the structure of the mold, it is extremely difficult to enable the plurality of floating cores to move in the different directions.\nPatent Document 3 discloses a method of manufacturing a pipe with a branch by molding using a floating core."}
{"text": "Many processing systems include multi-core processors in order to decrease the processing time to perform a task. A multi-core processor includes two or more cores coupled together. Hence, by dividing the operations of a processing task between the multiple cores, an improvement in the time required to complete the processing task can be achieved. However, this improvement is limited by the number of operations in the processing task that can be performed in parallel. For example, commutative operations, such as addition or multiplication of a series of numbers, can be performed in parallel to improve the processing time. However, other operations are serial in nature and cannot be substantially improved through parallel processing. For example, calculating a running total based on user input requires the processor to wait for the user input prior to performing each calculation. Thus, the operation is serial in nature and will not benefit substantially from parallel processing on a multi-core processor.\nParallel performance is also limited by the amount of communications overhead required to integrate the partial results that are computed on different cores. For example, in the case of adding n numbers together, each of p cores can add n/p of the numbers together in parallel. However, after this parallel step, there is a communications step in which the cores exchange the partial sums that have been computed. This is a communications overhead that does not occur when adding the numbers on a single core.\nIn order to benefit from the improvements offered by a multi-core processor, application code is often parallelized, where appropriate, such that each core is assigned a portion of the processing tasks that can be performed in parallel. However, current techniques for producing parallelized or concurrent application code suffer from various limitations. For example, one technique involves having a programmer write the application code using a parallel language designed for multi-core processing. However, manually producing the code using a parallel language is more prone to human error and places an additional burden on the programmer."}
{"text": "With the wide application of Chinese characters, learners of Chinese often fail to pronounce some word segments in a process of reading a Chinese character text aloud, especially some polyphonic word segments. Therefore, the learners of Chinese always come across obstacles of reading aloud. Even though pronunciations of a polyphonic word segment can be obtained by looking up a dictionary, much time is required to do so. In addition, because a vocabulary included in the dictionary is large, it is difficult for a user to grasp the most accurate pronunciation of the polyphonic word segment. This lowers the efficiency and effects in text reading, thereby affecting effects in Chinese learning."}
{"text": "This invention relates to a novel polypropylene composition having specific utility in the manufacture of oriented polypropylene films and to polypropylene films prepared therewith. More specifically, it relates to polypropylene compositions containing certain fatty acid amides not heretofore known to be useful as components in polypropylene compositions.\nIn the manufacture of polypropylene films, it is common practice to include, in the polymer formulation, additives referred to as slip agents. These additives migrate to the surface of the film and decrease the coefficient of friction between the film and the metal rollers over which it is passed during processing, thus facilitating the processing of the film. They also decrease the coefficient of friction between layers of the film when it is wound into rolls thereby facilitating unwinding of the rolls for further processing.\nThe slip additives employed in most commercial polypropylene films are relatively high molecular weight fatty acid amides. The most widely used fatty acid amides are erucamide, an unsaturated 22 carbon amide (13-docosenamide) and behenamide (docosanamide), the saturated analogue of erucamide. Both of these compounds are readily available, naturally occurring materials. They are normally provided as mixtures containing a small amount of other amides containing about 18 to 20 carbons.\nErucamide works quite well as a slip agent, but it is not favored by some producers and converters as it is relatively volatile and extra care is required during film manufacture and conversion to avoid having quantities of the additive escape from the film and plate out on processing equipment, thereby causing a clean-up problem. For this reason, some producers prefer to use behenamide as the slip agent. As slip agents, erucamide and behenamide appear to perform substantially equally well but behenamide is less volatile and therefore films containing behenamide are easier to handle during processing since the behenamide does not escape and plate out on the processing equipment.\nDespite the inconveniences encountered with erucamide-containing film, erucamide is the fatty acid amide slip agent of choice for those who wish to print the film using the well known rotogravure printing technique. In the rotogravure process, ink is applied to the print surface from a gravure cylinder containing the desired image as depressions on its surface. Ink is applied to the cylinder and the surface of the cylinder is wiped by a doctor blade to remove excess ink. The film to be printed is then contacted by the gravure cylinder and the image is thus transferred from the cylinder to the film. While the film is in contact with the cylinder, a small amount of erucamide is deposited on the raised portion of the cylinder. Since erucamide is soluble in the alcohol or ketone ink solvent, it is dissolved by the solvent in the next application of ink and is removed by the next swipe of the doctor blade. Thus, no build-up of erucamide occurs on the doctor blade or on the gravure cylinder.\nBehenamide, on the other hand, lacking the carbon-carbon unsaturation found in erucamide, is substantially less soluble in the alcohols and ketones found in printing inks. Thus, when, in the course of rotogravure printing, some of the behenamide is deposited on the printing cylinder, it is not dissolved in the printing ink and wiped off when the next application of ink and wiping by the doctor blade takes place. As a result, the behenamide is merely wiped off by the next swipe of the doctor blade and substantially all of it remains on the blade. After multiple swipes of the blade, a behenamide build-up forms on the doctor blade and on the printing cylinder, eventually reaching a point at which the blade is prevented from wiping the surface of the printing cylinder clean. This build-up causes streaks to form on the printed film, requiring an interruption of the job in order to clean the doctor blade.\nIt will be immediately apparent that a slip additive that does not cause the film processing problems associated with erucamide but which allows the film to be rotogravure printed without causing the problems associated with behenamide would be a welcome contribution to the polypropylene film art. In accordance with this invention, there has been discovered a slip additive that accomplishes this objective.\nIt is an objective of this invention to provide oriented polypropylene films that exhibit improved properties in rotogravure printing as compared to films presently known to the art. It is likewise an objective of this invention to provide film-forming compositions for preparing such films.\nIn accordance with this invention, it has been found that certain bis-fatty acid amides exhibit a combination of properties that makes them highly satisfactory as slip agents in oriented polypropylene films. In particular, this combination of properties make them highly satisfactory as slip agents in films that are to be printed in rotogravure printing operations.\nBriefly stated, the invention is a composition comprised of a polymer or copolymer of propylene and an N,Nxe2x80x2-bis-alkylene fatty acid amide having the general structural formula Rxe2x80x94COxe2x80x94NHxe2x80x94(CH2)nxe2x80x94NHxe2x80x94COxe2x80x94Rxe2x80x2 wherein R and Rxe2x80x2 are the same or different alkyl or alkenyl groups having about 15 to 21 carbon atoms xe2x80x94COxe2x80x94 is a carbonyl group and n is an integer from 2 to 4 said N,Nxe2x80x2-bis alkylene fatty acid amide being present in an amount from about 0.05 to about 0.5% by weight based on the combined weight of said amide and the polypropylene.\nIn another aspect, the invention contemplates an oriented film comprised of a polymer or copolymer of propylene and an N,Nxe2x80x2-bis-alkylene fatty amide having the general structural formula Rxe2x80x94COxe2x80x94NHxe2x80x94(CH2)nxe2x80x94NHxe2x80x94COxe2x80x94Rxe2x80x2 wherein R and Rxe2x80x2 are the same or different alkyl or alkenyl groups having about 15 to 21 carbon atoms, n is an integer from 2 to 4 and xe2x80x94COxe2x80x94 is a carbonyl group said N,Nxe2x80x2-bis alkylene fatty acid amide being present in an amount from about 0.05 to about 0.5% by weight based on the combined weight of said amide and the polypropylene.\nA preferred embodiment of the invention is a composite film comprised of a polypropylene core having, on at least one of its surfaces, a film layer comprised of a composition comprising polypropylene and up to about 0.5% of an N,Nxe2x80x2-bis-alkylene fatty amide having the general structural formula Rxe2x80x94COxe2x80x94NHxe2x80x94(CH2)nxe2x80x94NHxe2x80x94COxe2x80x94Rxe2x80x2 wherein R and Rxe2x80x2 are the same or different alkyl or alkenyl groups having about 15 to 21 carbon atoms, n is an integer from 2 to 4 and xe2x80x94COxe2x80x94 is a carbonyl group, said N,Nxe2x80x2-bis alkylene fatty acid amide being present in an amount from about 0.05 to about 0.5% by weight based on the combined weight of said amide and the polypropylene.\nAs used in the description of this invention and in the attached claims, the term xe2x80x9cfilmxe2x80x9d refers to a stand-alone film, i.e. a film of a sufficient thickness to have the strength and other characteristics required to have utility in packaging and other applications without having to be united with another film to support it. Stand-alone monolayer films are seldom seen in commerce. However, the compositions of the invention are suitable for use in such films.\nOriented polypropylene films of commerce are usually composite, i.e. multilayer, structures in which a core layer of a thickness sufficient to impart stand-alone properties to the overall structure carries one or more thin functional layers on its surface(s). The term xe2x80x9cfilmxe2x80x9d will also be used to refer to these thin functional layer films although they do not have sufficient thickness to stand alone.\nComposite polypropylene films are typically comprised of a homopolypropylene core having one or more functional layers on its surfaces. Functional layers can act, e.g., as heat or cold seal layers, as receiving layers for printing or other decorative material or as barrier coating receiving layers. In these cases, it is not required that the polymer in the surface layer be the same as that employed in the core, although it may be the same polymer formulated differently to serve the functional role that it is to serve in the completed film. Frequently, the surface layer will be of a different polymer that has been found to perform the desired function more satisfactorily than would the polymer employed in the core.\nWhen reference is made herein to polypropylene, it is intended to indicate a crystalline (isotactic) propylene homopolymer or a copolymer of propylene with ethylene or an xcex1-olefin having 2 to 5 carbon atoms in an amount insufficient to have a significant effect on the crystallinity of the polypropylene. Typically, the comonomer will be ethylene in an amount of 6% or less.\nSuitable polypropylenes are the commercially available isotactic polypropylenes having a melt flow rate between about 2 and about 10 dg/10 min at 230xc2x0 C. and 2.16 Kg. load and a DSC melting point of about 160 to 166xc2x0 C. One polypropylene that can be used is the isotactic homopolymer having a melt flow rate of about 3.5 dg/1 0 min at 230xc2x0 C. and 2.16 Kg. load, available from Aristech Chemical Corporation, Pittsburgh, Pa. Suitable polypropylenes are also available from Montell, Inc. Wilmington, Del., Exxon Chemical Company, Baytown, Tex. and Fina Oil and Chemical Co., Deer Park, Tex.\nSince different layers are formulated differently, it is possible that different slip agents may be used in the different layers or that some layers contain no slip agent. Thus, it is possible and desirable, in cases where problems are encountered when using conventional slip agents, e.g. in the rotogravure printing process as described above, to limit the use of the bis-amides of the invention to the functional layer where its presence yields the desired advantage. In such cases, where in the past the conventional slip agents have been added to the core layer and allowed or forced to migrate into the functional layers, it is found that the presence of the bis amides almost completely prevents the migration of the conventional slip agent, e.g. erucamide or behenamide from the core to the functional layer where it can cause a problem.\nThe slip additives employed in the compositions and oriented films according to this invention are N,Nxe2x80x2-bis- fatty acid amides having the general structural formula Rxe2x80x94COxe2x80x94NHxe2x80x94(CH2)nxe2x80x94NHxe2x80x94COxe2x80x94Rxe2x80x2 wherein R and Rxe2x80x2 are the same or different alkyl or alkenyl groups having about 15 to 21 carbon atoms, n an integer from 2 to 4 and xe2x80x94COxe2x80x94 is a carbonyl group. The compounds employed in the compositions of this invention are characterized by having a relatively low degree of volatility and a relatively high degree of solubility in alcohols, ketones and esters.\nCompounds meeting the above definition and description include, by way of example, N,Nxe2x80x2bis-ethylene stearamide, N,Nxe2x80x2-bis ethylene oleamide, N,Nxe2x80x2-bis ethylene behenamide, N,Nxe2x80x2-bis ethylene erucamide, N,Nxe2x80x2-bis propylene stearamide, N,Nxe2x80x2bis butylene oleamide.\nThe N,Nxe2x80x2-bis-alkylene fatty acid amide is present in the compositions of the invention in a concentration of about 0.05 to about 0.5% by weight based on the total weight of the composition. Preferably, the concentration of the slip additive is between about 0.05 and about 0.25% and most preferably between about 0.5 and about 0.2%.\nThe slip agents of the prior art such as erucamide or behenamide are sufficiently volatile that a measurable portion of them always migrates from the interior of the film and forms a layer of slip agent on the film\"\"s surface, whether the film is a composite or a monolayer. The amount present on the surface can be measured by a technique referred to as Attenuated Total Reflectance Infra-Red Spectroscopy (ATRIR). When the ATRIR spectrum of a film sample stored at room temperature is compared to that of a film stored for 16 hours at 55xc2x0 C., an increase in the amount of a prior art slip agent is detected. (Such heat aging is commonly applied to simulate film aging conditions.) The amount of the slip agents of this invention that migrates to the surface of a film cannot be measured by ATRIR. Their presence on the film surface is demonstrated only by their efficacy as slip agents.\nThe presence of conventional slip agents on the surface of a film can also be detected by rinsing the film surface with a solvent and testing for the slip agent by High Pressure Liquid Chromatography. Similarly, no accumulation of the slip agent of the invention can measured on the film surface by this technique, whereas it is readily detected in the solvent employed to wash the surface of a prior art film.\nIn prior art practice, slip agent is frequently incorporated into the core layer of composite films, which are then heat treated to force it to migrate to the surface layers. In other cases, if the slip agent is present only in a surface layer, it tends to migrate into the core layer although, as a rule, it is not needed there. In either case, this migratory tendency results in slip agent being present in locations where it is not needed. Moreover, it also requires the use of greater quantities of slip agent than are needed in order to assure that sufficient quantities are present in the area where the real need exists.\nIt has been found that the slip agents of this invention exhibit a very low tendency to migrate between layers or even within a layer of the film. Thus, a slip agent incorporated in a surface layer does not migrate into the core layer, nor does it migrate and accumulate on the surface of a layer as do the slip agents of the prior art. When tested by the ATRIR test or the surface rinse procedure mentioned above no accumulation of the slip agent is detected on the surfaces. The low migratory character of these slip agents causes them to remain in the layer into which they are incorporated.\nDespite the absence of measurable amounts of the slip agent on the film surface, the slip agents of the invention are substantially equal to those of the prior art in decreasing the coefficient of friction (COF) of the film surface. Thus, the COF measured on the surface of a prior art film containing erucamide as the slip agent is about 0.15 to 0.35. That measured on a prior art film containing behenamide is about 0.15 to 0.35. The COF measured on the surface of a film of this invention is generally between about 0.15 and about 0.45.\nThe low migratory character of the slip agents of the invention also leads to other favorable effects on film processing. For example, the problem of vaporization encountered with erucamide and to a lesser extent with behenamide is eliminated. Also, as will be apparent from the above, in the preparation of composite films, it is not always necessary to have a slip agent in the core layer. Since the slip agent of the invention will not migrate into the core layer, one can put only the amount actually needed to give the desired slip properties in the surface layer and it will not migrate into the core layer and reduce the level in the surface layer below that required to give the surface of the film the proper slip quality.\nAs indicated hereinabove, the compositions of the invention have particular utility in the manufacture of rotogravure printable oriented polypropylene film thanks to the unique combination of good alcohol, ketone and ester solubility and low volatility exhibited by the N,Nxe2x80x2-bis-alkylene fatty acid amides. Due to their low degree of volatility, these compounds do not vaporize from the surface of the film during manufacturing and conversion. Due to their relatively high degree of solubility in the solvents that are normally employed as solvents for printing inks, whatever small amounts these compounds that might be on the film surface does not form a build-up on the doctor blade during rotogravure printing. Films prepared using these compositions thus combine the favorable properties of conventional films containing behenamide during the manufacture and conversion of the film with the favorable properties of those containing erucamide for use in rotogravure printing.\nIn yet another aspect, the slip agents employed in the films of the invention are found to be improved candidates for use as cold seal films. Cold seal films are prepared by treating the films with a cold sealable adhesive such as a rubber cement and many applications for printed films are applications in which cold seals are preferred over more conventional heat seals. In prior art usage, particularly when erucamide is employed as the slip agent, it is found that the high level of the slip agent on the film surface weakens or even, in extreme cases, destroys cold sealability. Using the slip agents of this invention at the levels specified, very little or no harmful effect on cold seal is seen.\nIn addition to the slip agents of this invention, conventional additives, in conventional amounts, can also be included in the films. Suitable conventional additives include, by way of example, antioxidants, pigments, orientation stress modifiers, flame retardants, antistatic agents, antiblocking agents, and antifoggants. Any such additives that do not interfere with the functioning of the slip agent can be present.\nIn any embodiment of the films of the invention, the film can be either clear and transparent or it can be opaque. Typically, polypropylene films are rendered opaque by loading the core layer with a void forming opacifier such as calcium carbonate. Such opacifiers cause opacity by forming microvoids in the polymer matrix during the drawing operation. Other types of opacifiers which cause opacity simply by increasing the optical density without voiding can also be employed. Typical of such an opacifier is titanium dioxide.\nThe total film thickness, for either the single or multiple layer films, is preferably in the range of about 0.25 to about 1.5 mil, i.e. about 25 to about 150 gauge. In the multilayer embodiment, the core layer preferably has a thickness of about 23 to about 40 gauge, while the other layers preferably each have a thickness of about 2 to about 10 gauge.\nTo prepare a printable surface, the surface of the film that is to be printed is subjected to an oxidation treatment to impart a degree of polarity to the surface, which ensures good adhesion of the printed information to the film. Suitable oxidation treatments include, by way of example, corona discharge, flame and acid etching. The preferred treatment is with corona discharge or with flame. The use of flame and corona, seriatim, is also known and frequently employed.\nThe oxidation treatment is preferably carried out to a degree sufficient to create a surface tension on the film surface equal to about 35 to about 60 dynes/cm. A preferred surface tension is about 35 to about 40 dynes/cm.\nAs stated hereinabove, composite films can have, in addition to the printable surface, other functional layers on the surface opposite the print layer. If the layer on the side opposite the print layer is a barrier coating receiving layer, that layer will also be subjected to an oxidation treatment to approximately the same degree as is applied to the print layer.\nWhile the invention has been described with emphasis on its utility in rotogravure printing, it will be apparent that other types of films, including functional layers, can also benefit from the low volatility and the lower degree of migratory tendency of the slip agents of the invention. One type of such film is metallizable film where the greater migratory tendency and volatility of conventional slip agents tends to result in a greater amount of slip agent on the film surface, which can sometimes interfere with bonding of the metal coating to the film surface. The same effect is also noted with cold seal coatings applied to the film surface.\nPolypropylene films according to the invention can be prepared using film forming and orientation methods and techniques well known to the plastic film art. Thus the films can be prepared using the tubular or bubble process wherein a tube of film is extruded through a circular die and is drawn and inflated simultaneously to effect biaxial orientation. Likewise, they can be prepared by the tenter process in which a flat sheet of the film is drawn from the die at a multiple of the lineal rate of extrusion to effect longitudinal orientation and is thereafter drawn laterally by a preselected multiple of its original width to effect lateral orientation. As a rule, films prepared by the bubble process are drawn uniformly 6xc3x97 by 6xc3x97 or 7xc3x97 by 7xc3x97. When the tenter process is employed, the draw is usually carried out to 5xc3x97 in the machine direction and 10xc3x97 in the cross direction. These draw ratios are also suitable for the films of the invention, although other ratios can be employed if desired."}
{"text": "1. Field of the Invention\nThe present invention relates to the providing of anesthesia during surgery, and particularly to devices and methods for coordinating and organizing the numerous patient tubes, catheters, wires, monitoring lines and other objects used by an anesthesiologist during surgery, particularly major surgery. The invention is designed to travel with the patient for use by other health care practitioners following surgery, particularly in the intensive care unit and in the recovery room.\n2. Description of the Prior Art\nWhen a patient arrives at the operating room (OR) for major surgery (such as, for example, heart, major chest, abdominal, renal or liver surgery), the anesthesiologist will ordinarily attach a host of invasive tubes, lines, catheters, wires and other devices to the patient that are used throughout the surgery and often afterwards. These include such things as one or two intravenous peripheral lines that are attached to an upper extremity of the patient; an arterial line that is attached to an arm or the groin of the patient; one or two central lines that are attached to the neck or upper chest of the patient; a cordis line that is attached to the neck or upper chest of the patient; and a Swan-Ganz catheter attached to the neck or upper chest of the patient.\nDuring and after surgery, the patient will need drugs (such as, for example, vasopressor and other supportive drugs) some of which are administered via dripping methods, and others via injection into tubes. The patient may also require blood, plasma and other fluids. These materials and drugs may be supplied directly to the patient through the central line, or through one of the lines extending through the Swan-Ganz catheter. Depending upon which Swan-Ganz catheter is used, it may have several ports which will be connected to different tubing, monitoring cables and pressure transducers. One of these ports may also be connected to a drip line that holds several stopcocks which could receive different medications from the dripping machines to be given to the patient.\nAll of these lines, tubes, and catheters come from different parts of the patient body and lead around the head of the patient where the anesthesiologist is stationed during surgery. Here, the multitude of lines and tubes are connected to monitors, screens, drip lines, injection ports and other devices that the anesthesiologist must keep control of during the hours of surgery. It is easy to confuse these numerous lines and tubes. This can be dangerous in the event that an emergency situation arises during surgery where the anesthesiologist is required to quickly provide a drug or other medication to the patient. Delay by the anesthesiologist while trying to locate the correct tube to administer the drug, or inadvertent administration of a drug to the wrong site could endanger the life of the patient and affect the outcome of the surgery.\nDuring lengthy surgery, a first anesthesiologist may be relieved by a second anesthesiologist who must step in and figure out how the tubes, lines and wires have been set up by the first anesthesiologist. Since there is no established or standard way to set up these tubes, lines and wires, the second anesthesiologist must figure out how the first anesthesiologist set everything up. If the transition from one health care practitioner to another is not smooth, there is another dangerous potential for delay or error by the second anesthesiologist in administering a needed drug during an emergency situation. Similarly, the maze of unorganized tubes, lines, tape, labels and devices make it difficult to teach student doctors (residents and fellows).\nIt is also very easy for the tubes and lines to become tangled, and unless they are well secured and supported, it is possible that they can be pulled out, or become dislodged or disconnected—sometimes without the knowledge of the anesthesiologist—placing the patient in a threatening and risky situation. This is especially true for the Swan-Ganz catheter which can cause serious or potentially fatal results if dislodged.\nOnce surgery is completed, the patient must be transported from the operating room to an intensive care unit or elsewhere. Many of the monitors, drip lines and other tubes must remain inserted into the patient for hours or days after surgery. The simple act of transporting the patient down the hall from an operating room to an intensive care unit requires bringing all of these wires, tubes and lines along too—as well as the machines, drip bags/bottles, pacemaker, and other devices that they are connected to. There is a serious risk of pulling out a tube, line or catheter as all of these things are manipulated during transport.\nSome of the devices must be disconnected while the patient is traveling, and reconnected when the patient arrives at the intensive care unit. Once the patient arrives at the intensive care unit for recovery, the nurses and other health care providers must again untangle and sort out all of the tubes and lines, and reconnect the lines and wires to their associated monitors and other devices. Some of the monitors provide critical data regarding the condition of the patient, such that any unnecessary delay in reconnecting could pose serious risks to the patient.\nThere are some organizational tools in the prior art that begin to address these situations, including the invention disclosed in U.S. Pat. No. 4,988,062. This patent discloses a multi-part pad with pivotally attached wings, the wings having manifold devices attached at their distal ends for receiving and separating the tubes and lines coming out of the patient. However, this device is not designed to travel with the patient, and once the tubes are placed in the manifolds, the wings may still be pivoted which could cause the tubes to be pulled out or disconnected, having potentially disastrous results. Another device disclosed in U.S. Pat. No. 5,624,403 is a removable pad that is designed to be placed across the torso of a patient, the upwardly facing surface of the pad being composed of loop pile material (such as Velcro®) for receiving corresponding strips that hold down the tubes and lines leading from the patient. However, this device is designed to stretch across the torso of the patient, and therefore cannot be used during abdominal surgery, nor can it be used by an anesthesiologist who is stationed above the head of the patient. There are also elaborate trays such as those disclosed in U.S. Pat. Nos. 4,720,881, 5,334,186 and 5,435,448. None of these prior art inventions provides an organizational guide for the numerous tubes, lines, catheters and wires used by an anesthesiologist during major surgery; none of them are designed to travel with the patient from the operating room to intensive care or elsewhere for use by other health care providers; and none of them serves as an educational tool or reference for the new anesthesia student, resident, fellow or new graduate anesthesiologist to learn about Swan-Ganz catheter usage, pacemaker modes, doses, and administration of commonly used drugs during surgery.\nIt is therefore desirable to provide devices and methods for coordinating and organizing the numerous patient tubes, catheters, wires, monitoring lines as well as other objects and machines (e.g., pacemaker, pressure monitors) used by an anesthesiologist during major surgery; to provide standardization for the positions of these various tubes, wires and lines so that the anesthesiologist, his/her replacement and other health care providers (particularly residents and fellows) can know instantly where everything is; to provide coordination and organization devices that prevent inadvertent disengagement of tubes, lines and wires leading from the patient during and after surgery; and to provide coordination and organization devices that can travel with the patient and remain with the patient during recovery for hours or days after surgery."}
{"text": "1. Field of the Invention\nThe present invention is related to electronic memories for data processors and signal processors.\n2. Description of the Prior Art\nMultitudes of different types of memories are used in the prior art including magnetic memories, integrated circuit memories, and optical memories. The most pertinent prior art is integrated circuit read only memories and integrated circuit charge coupled device (CCD) memories. Integrated circuit read only memories are significantly more expensive than charge coupled device memories, primarily because such read only memories require more complex monolithic processes than are required for CCD memories and because such read only memories are random access memories thereby requiring more complex accessing circuitry and multitudes of parallel interconnections than are required for CCD memories. CCD memories are lower in cost than conventional read only memories, but CCD memories are volatile where all stored information is lost when power is turned off. This volatility characteristic is verified in the \"Engineer's newsletter\"; Electronics Magazine; Mar. 31, 1977; page 121; last sentence.\nTherefore, the prior art has been constrained to expensive memories such as magnetic memories and integrated circuit read only memories when non-volatile capability is required.\nThe prior art is further defined in the art-of-record of the related applications including U.S. Pat. No. 3,356,989 to Autry; U.S. Pat. No. 3,613,777 to Quay; U.S. Pat. No. 3,643,106 to Berwin; U.S. Pat. No. 3,826,926 to White; U.S. Pat. No. 3,852,745 to LeBail; U.S. Pat. No. 3,873,958 to Whitehouse; U.S. Pat. No. 3,876,989 to Bankowski; U.S. Pat. No. 3,889,245 to Gosney; U.S. Pat. No. 3,775,738 to Quay; U.S. Pat. No. 3,787,852 to Puckette; U.S. Pat. No. 3,891,977 to Amelio; U.S. Pat. No. 3,895,342 to Mallet; U.S. Pat. No. 3,909,806 to Uchida; U.S. Pat. No. 3,914,748 to Barton; and U.S. Pat. No. 3,942,034 to Buss and including the article by Dennard, IBM Technical Disclosure Bulletin, Vol 14, No 12, May 1972, pages 3791-3792; the article by Altman, Electronics Magazine, Feb. 28, 1972, pages 62-71; and the article by Baertsch, Electronics Magazine, Dec. 6, 1971, pages 86-91 which references are all incorporated herein by reference."}
{"text": "The present invention relates to a method for aligning a pair of planar templates on opposite sides of a planar board, and apparatus for practicing this method.\nDouble-sided circuit boards having aligned circuit patterns on opposite board sides are commonly used in the electronics industry. In the usual case, such a board is formed by placing circuit pattern layouts at mutually aligned positions on opposite sides of the board. A photo-etching process is used to remove copper selectively from plating on the two board sides, so as to leave the desired circuit pattern on opposite sides of the board.\nThe present invention contemplates a method and apparatus for aligning a pair of templates on opposite sides of a board. The method of the present invention includes cradling a board and a template placed against one side thereof, with such cradling producing alignment in the template and the board. Where the templates and board have matching, intersecting edges, the template and board are supported, in their cradled position, at such edges. The aligned board and template are flipped to expose the board's other side. The second template is placed against the exposed side of the board, and the two templates and board are cradled, again to produce alignment between the board and the second-mentioned template.\nThe apparatus of the invention includes structure defining a pair of surfaces, each surface being adapted to promote alignment between a board and a template placed against one side of the board. The surfaces are mounted for shifting toward and away from a position wherein a board and template placed against one are flipped onto the other to expose the other side of the board.\nA preferred embodiment of the invention includes an open-book structure having a pair of interior surfaces which intersect along an inclined edge having raised and lowered ends. The structure pivots about an axis substantially paralleling this edge. Projecting from each of the surfaces in the structure, adjacent the lower end of the inclined edge, is a pin which cooperates with the two surfaces to promote alignment of a board and template placed against one of the surfaces.\nA general object of the present invention is to provide a simple, efficient method for aligning a pair of substantially planar templates on opposite sides of a planar board, and apparatus for practicing the method.\nA more specific object of the invention is to provide such a method wherein alignment between a board and a template placed against one side of the board is produced by cradling the board and template."}
{"text": "1. Field of the Invention\nThe invention relates generally to navigation systems and industrial control systems and more specifically to combinations of navigation and industrial controls that provide centimeter accuracy in the positioning and movement of operating machinery and robots.\n2. Description of the Prior Art\nWhen originally conceived, the global positioning system (GPS) that was made operational by the United States Government was not foreseen as being able to provide centimeter-level position accuracies. Such accuracies are now commonplace.\nExtremely accurate GPS receivers depend on phase measurements of the radio carriers that they receive from various orbiting GPS satellites. Less accurate GPS receivers simply develop the pseudoranges to each visible satellite based on the time codes being sent. Within the granularity of a single time code, the carrier phase can be measured and used to compute range distance as a multiple of the fundamental carrier wavelength. GPS signal transmissions are on two synchronous, but separate, carrier frequencies \"L1\" and \"L2\", with wavelengths of nineteen and twenty-four centimeters, respectively. Thus within nineteen or twenty-four centimeters, the phase of the GPS carrier signal will change 360.degree..\nHowever, the number of whole cycle (360.degree.) carrier phase shifts between a particular GPS satellite and the GPS receiver must be resolved. At the receiver, every cycle will appear the same. Therefore there is an \"integer ambiguity\". The computational resolution of the integer ambiguity has traditionally been an intensive arithmetic problem for the computers used to implement GPS receivers. The traditional approaches to such integer ambiguity resolution have prevented on-the-fly solution measurement updates for moving GPS receivers with centimeter accurate outputs. Very often, such highly accurate GPS receivers have required long periods of motionlessness to produce a first and subsequent position fix.\nThere are numerous prior art methods for resolving integer ambiguities. These include integer searches, multiple antennas, multiple GPS observables, motion-based approaches, and external aiding. Search techniques often require significant computation time and are vulnerable to erroneous solutions or when only a few satellites are visible. More antennas can improve reliability considerably. If carried to an extreme, a phased array of antennas results whereby the integers are completely unambiguous and searching is unnecessary. But for economy, the minimum number of antennas required to quickly and unambiguously resolve the integers, even in the presence of noise, is preferred.\nOne method for integer resolution is to make use of the other observables that modulate a GPS timer. The pseudo-random code can be used as a coarse indicator of differential range, although it is very susceptible to multipath problems. Differentiating the L1 and L2 carriers provides a longer effective wavelength, and reduces the search space. However, dual frequency receivers are expensive because they are more complicated. Motion-based integer resolution methods make use of additional information provided by platform or satellite motion. But such motion may not always be present when it is needed.\nAnother prior art method and apparatus for precision attitude determination and kinematic positioning is described by Hatch, in U.S. Pat. No. 4,963,889, comprises the steps of determining the approximate initial relative position of a secondary antenna that is freely movable with respect to a reference antenna; making carrier phase measurements based on the reception of \"N\" number of satellites, where N is the minimum number of satellites needed to compute the relative position of the secondary antenna; deriving from the carrier phase measurements an initial set of potential solutions for the relative position, wherein the initial set of potential solutions all fall within a region of uncertainty defined by a sphere having a radius equal to the maximum distance between the two antennas, and wherein multiple potential solutions arise because of whole-cycle ambiguity of the carrier signal; making redundant carrier phase measurements based on the reception of a carrier signal from an additional satellite (N+1); and eliminating false solutions from the initial set of potential solutions, based on a comparison of the redundant carrier phase measurements with the initial set of potential solutions, to reduce number of potential solutions to close to one, whereby the number of potential solutions is not increased by use of the redundant carrier phase measurements.\n\"Deriving from the carrier phase measurements an initial set of potential solutions\" means that the initial set are derived from just two satellites. The rest of the Hatch specification explains why N is exactly two in the case of attitude determination. Planar intersections of wave fronts are formed from the two satellites, thus obtaining a collection of parallel lines. The intersection points of these lines and a baseline sphere are determined, producing the initial set of potential solutions. For example, as explained at column 12, line 35 of Hatch, there are 188 points or potential solutions in the initial set. In the Hatch method, \"eliminating false solutions from the initial set of potential solutions,\" means eliminating 187 of those 188 points. The idea of potential solutions refers to the initial set of 188 points.\nHatch forms an initial collection of around 188 potential solutions using just two satellites, and then uses phase measurements of the remaining satellites to whittle away at that small initial collection, leaving only one candidate solution if phase measurements are accurate enough. Hatch avoids having to deal with large numbers of integer combinations.\nDonald Knight, in U.S. patent application, titled \"METHOD AND APPARATUS FOR MAXIMUM LIKELIHOOD ESTIMATION DIRECT INTEGER SEARCH IN DIFFERENTIAL CARRIER PHASE ATTITUDE DETERMINATION SYSTEMS\", Ser. No. 07/976,174, filed Nov. 13, 1992, and incorporated herein by reference, describes a method of reducing the mathematical intensity of integer ambiguity resolution.\nAside from making carrier phase measurements and determining a region of uncertainty, Hatch's method is essentially a two-step method, which is a way to avoid the problem of large numbers. The essence of the Knight patent is to do the opposite, plunging into a sea of possible solutions, even if there are trillions of them, and finding the single, best solution out of the entire product space of possibilities.\nIf ten different whole cycle values can be added to the first difference phase measurements from each of four satellites, then Knight defines 10.sup.4 or 10,000 possible acceptable solutions. Knight then finds the one-and-only solution out of 10,000 that matches the measured phase data better than all 9,999 other combinations. If eight satellites are available (which does happen), then Knight finds the single best solution out of 10.sup.8 or 100 million possibilities. If two antenna pairs are present, then the number becomes (10.sup.8).sup.2 =10.sup.16. Thus whereas Hatch uses a two-step method of avoiding large numbers, Knight provides a practical method for finding the single best solution in a solution space with thousands or millions of possibilities. Hatch even suggests that the Knight approach is incorrect. For example, at column 8, Line 63, Hatch states: \"if ten different whole cycle values can be added to the first differences from each of the four satellites, then one might expect 10.sup.3 or 10,000 possible solutions within the uncertainty region . . . . Indeed, some authors have defined the number of possible acceptable solutions in exactly this fashion.\"\nKnight describes a method that is more computationally intensive than Hatch's, and has proven to be reliable and workable in commercially-viable applications. Hatch's method avoids the problem of large numbers, and does best under conditions in which the problem of large numbers does not exist, e.g., under conditions where the myriad of possible solutions collapse into a handful. In the real world, conditions are usually not optimum, which would require the corresponding equipment that uses Hatch's method to be more complex, in order to achieve the same results as the present invention's.\nHatch's two-step method includes preliminary steps involving making carrier phase measurements, and then numbering them one through N, followed by an (N+1).sup.th measurement. It could be expected that any similar apparatus will make carrier phase measurements, and the basic idea of interferometry is conventional, regardless of how the carrier phase measurements are numbered. Such preliminary steps define which carrier phase measurements are involved in generating an initial set of potential solutions, and which carrier phase measurement is defined to be the redundant measurement for eliminating potential solutions. As applied to attitude determination, the step of determining the approximate relative position of the secondary antenna equates to measuring the distance between the antennas, and noticing that the secondary antenna has to lie on a sphere centered at the main antenna.\nHatch's method involving a \"plurality N of satellites\" is not the same as a \"plurality of satellites\". As applied to attitude determination, N is exactly two. As applied to kinematic surveying, N is exactly four. Hatch makes this explicit in FIG. 7, blocks 58 and 62, and states at column 12, line 3, \"The first computations to be performed in an attitude determination application are to compute the entire set of potential solutions based on the carrier phase measurements from two satellites. One of the significant features of the method described is the use of only two of the satellites to determine the set of potential solutions, rather than more than two satellites. Hatch explains that only a plurality of N=2 satellites are required\" . . . in the attitude determination application, . . . there is! only a two dimensional uncertainty region corresponding to the surface . . . . Thus only two satellites are needed to define all the potential solutions that can exist on the two dimensional uncertainty region . . . . The attitude detection application is a two dimensional uncertainty problem, requiring two satellites to define all of the potential solutions . . . . The final broad step in the method of the invention is eliminating false potential solutions based on redundant information from additional satellites, i.e. additional to the minimum two . . . satellites required to define all of the potential solutions.\nThe phrases \"initial set\", \"initial set of potential solutions\" and \"possible solution\", refer specifically to the result obtainable by applying the first step of Hatch's two-step method. This \"initial set of potential solutions\", is to be generated from the received carrier signals of two and only two satellites, since the teaching is that \"only two satellites are required to define all of the potential solutions.\"\nThe notion of \"potential solutions\" is thus restricted to 188 points, or to some such number in a similar example, but certainly not 10,000 solutions. It is also very clear that \"eliminating false solutions from the initial set of potential solutions\" means eliminating 187 points out of 188 points.\nHatch contends that it is not necessary to consider 10,000 possible solutions, because \"the uncertainty region is two dimensional\" and the measurements from any two satellites span the region of uncertainty. Using two satellites results in a collection of only 188 possibilities. Hatch contends that the correct solution is one of those 188 points, and that those 188 points are the only possibilities, that those other myriad of possibilities do not exist.\nThe \"initial set of potential solutions\" of Hatch cannot be extended to cover an entire product space of 10,000 solutions, because the essence of the two-step method is avoiding the problem of large numbers. Hatch proposes reducing the dimensionally of the solution space to the bare minimum and declares that the right answer has to lie within that reduced solution space, and that only the first measurements corresponding in number to the bare dimensionally are required to determine all the possible solutions.\nKnight does not include an initial set. The method of Knight is to search for the right answer somewhere among 10,000 possibilities and finds the single best match out of such possibilities. There is no reason to stop and look at 188 possibilities when the remainder of 10,000 possibilities have to be searched. Furthermore, those 188 possibilities defined by Hatch do not resemble anything in Knight. Knight does not involve deriving from the carrier phase measurements an initial set of potential solutions for the relative position. Furthermore, Hatch generates those 188 points using a nonlinear operation on the carrier phase measurements of two satellites, something that Knight does not.\nHatch eliminates false solutions from the initial set of potential solutions, whereby the number of potential solutions is not increased by use of the redundant carrier phase measurements. Knight does not deal with Hatch's \"initial set\", whether to create it or eliminate from it. Most of the time, in a commercially-viable apparatus, the right solution will not be found in such an initial set.\nIf carrier phases could be measured perfectly, without radio reflections, reception noise, non-common antenna phase center motion, temperature sensitive drifts and other corruption of the measured phase data, then the attitude determination problem really would be two-dimensional. The right solution for position of the secondary antenna might then be found within Hatch's initial set of 188 potential solutions. Determining the right solution would be a simple matter of eliminating 187 false solutions by comparison with the third carrier phase measurement. As long as phase measurement errors are small enough, the two-step method based on near perfect measurements is satisfactory. Though there is some small perturbation in computed results, the right answer comes out in the end most of the time. Phase measurement errors of about two degrees RMS are usually tolerable.\nOnly one whole cycle value results in a surface that passes closer to a specific potential solution than any other whole cycle value. The problem is that phase measurements contain error, and geometric considerations amplify this error when measurements are processed in a computing apparatus. There is only one surface that passes closer to a specific potential solution than any other, but unfortunately, the surface that passes closest is often the wrong surface. Another surface farther away may be the right one. Not only can the surface be in the wrong place due to measurement error, Hatch's potential solution may be in the wrong place as well, due to error effects that are worse for the potential solution than the surface.\nThe supposed potential solution is actually a computational result that is often badly perturbed by phase measurement error combined with geometrical dilution of precision. A scissors effect can amplify measurement errors that distort the planes. Subsequently, the formation of intersections of lines with a sphere introduces a nonlinear operation on noisy data which greatly amplifies some of the errors. Hatch actually chooses the two satellites closest together in order to form an initial set with fewer than 188 points, whereas considerations of error sensitivity would dictate the opposite.\nHatch teaches that only one of the whole cycle values passes closer to a specific potential solution than any other. The two-step Hatch method prescribes that only the closest whole cycle value (from the additional satellite data) is selected for each of the potential solutions whereby the number of potential solutions is not increased. At some point, as phase measurement errors are increased, corresponding to less favorable operating conditions, the prescription begins picking the wrong whole cycle values a significant fraction of the time. Using Hatch's method, there is no way to pursue any option but the closest, and processing a phase measurement with the wrong whole cycle value added to it further accelerates the deterioration of computed results.\nBy comparison, all of the whole cycle values are tried as possibilities in the Knight method. An integer combination that looks very poor initially becomes more attractive as the solution progresses, eventually finishing as the best possible solution out of 10,000 or more. Every possible integer combination is pursued until, taken as a whole, it can be decided that the combination under development is less likely than another combination that has been completely finished and weighed with all satellite data included. The Knight method does not minimize the solution space dimensionally to the bare minimum and then restricts attention to an \"initial set of solutions\" obtained with only the data of the first measurements corresponding to that bare dimensionally. Most of the time, under real world conditions, the right answer simply is not to be found. The job of finding the right solution cannot be accomplished by avoiding the problem of large numbers.\nBenjamin W. Remondi, reported the status of differential GPS (DGPS) accuracy and on-the-fly (OTF) kinematic GPS (KGPS) in a paper titled, \"On-The-Fly Kinematic GPS Results Using Full-Wavelength Dual-Frequency Carrier Ranges\".\nOn-the-fly real-time kinematic (OTF-RTK) positioning systems are useful in earth moving machinery control and guidance, robotic applications and the placement of structures, such as bridge sections. For example, in Japan, a particular construction project in Tokyo Bay requires the highly accurate placement of bridge pilings and sections in the bay from the decks of ships that are rolling and drifting with the sea. Highly accurate OTF-RTK control systems are also useful in guiding unattended bulldozers in the path of volcano lava flows to redirect the flow of lava in hot, hostile environments.\nPhilip M. Clegg describes an automated earth grading system in U.S. Pat. No. 4,807,131. A powered grading machine is combined with a laser level to control the blade of the grader. The control system described includes an elevation control, a location positioning device, a grader-blade tilt detector, a digital terrain model, a digital computer for terrain comparisons, a blade servo controller and a visual display unit. The automated moving of earth for a wide-range of engineering applications is described. The system computes the location and elevation of a grader-blade and then compares this with a corresponding design height/grade specified in a computerized digital terrain model. If the design height is lower than the current height of the blade, then a cut is required, and the blade is lowered automatically or manually by an operator.\nOf particular interest is the use of satellite-based positioning systems, such as the global positioning system (GPS). Clegg mentions the use of celestial satellite positioning systems for the location of the grader. Many engineering applications require precise horizontal and vertical positioning. Centimeter-level GPS is typically achieved by using carrier phase and solving an integer lane ambiguity problem. Clegg does not teach the specifics of GPS nor how it might be used with his device.\nThe carrier phase ambiguity resolution problem has been the focus of a great deal of attention from research communities. In essence, the integer number of carrier cycles that existed between GPS receiver and satellite at the time that the signal was acquired must be determined. Direct range measurements, combined with the satellite geometry, allow the correct integer carrier phase ambiguities to be determined for a plurality of satellites tracked at two or more sites. Ambiguity resolution techniques that only use GPS measurements have become computationally efficient. The use of additional sensors such as a laser level, electronic distance meter, a compass, a tape, etc., provide valuable constraints that limit the number of possible integer ambiguities that need to be considered in a search for the correct set. For example see, INTEGRATED TERRESTRIAL SURVEY AND SATELLITE POSITIONING SYSTEM, filed Jul. 1, 1993, U.S. patent application Ser. No. 08/086,665."}
{"text": "The present invention relates to a poker-type card game and entertainment apparatus for playing the game.\nGambling is an exceedingly common form of entertainment to a great number of people. Players tend to like games that provide a chance of \"winning big\", are exciting to play, are simple to play, utilize a reasonable amount of thinking ability, and have a high payout ratio. For the casino, desirable games are fast so that the game can be played repeatedly in a very short period of time with each play only taking a matter of a few seconds.\nCasinos have long utilized machines to be operated by the players obviating the need for dealers. An exceedingly common form of such machines have been referred to as \"slot machines\". Slot machines do not require the need of a dealer and permit the player to play at his/her pace at the particular machine.\nAt one time, slot machines utilized only a plurality of rotating tumblers with indicia, such as cherries, lemons, oranges and other fruit, being represented on the tumblers. When the same type of fruit is in alignment at a horizontal line, the slot machines paid off a winning amount.\nWithin recent years, the slot machine has been expanded to be utilized with other types of games. For example, machines have been utilized to play blackjack and poker. Such machines have been readily accepted and provide a way of playing blackjack and poker at wagering levels which are substantially less than that would be required at a table for such a game where a dealer is utilized."}
{"text": "In 1962, Hanna-Barbera Production, Inc. produced a cartoon titled “The Jetsons” that featured the Jetson family living in a utopian future having incredible conveniences based on communication devices, large and small displays, and even robots. Everything they wanted required only a voice command or a push of a button.\nIn 1966, a television show titled “Star Trek” debuted that followed the interstellar adventures of James T. Kirk and the crew of the Enterprise. Star Trek is known for its influence on the world outside of science fiction for inspiring inventions such as desktop and tablet computers, wireless head phones, biometrics, cell phones, and even the automated sliding door.\nAn epic science fiction film released in 1968 titled “2001: A Space Odyssey” follows a voyage of a U.S. spacecraft to Jupiter. The film had tablet computers like Star Trek but also had two-way video conferencing, a suitcase phone, a voice controlled chessboard, and a sentient computer named HAL having artificial intelligence.\nForty-three years later, Apple® introduced the iPhone 4S®, which includes a computer program called Siri® that works as an intelligent personal assistant and knowledge navigator. Siri uses a natural language user interface to answer questions, make recommendations, and perform actions. Siri adapts to a user's individual language usage and preferences and returns individualized results.\nMany vehicles now offer various features activated by voice commands such as making telephone calls, selecting music, and interfacing with car navigation systems to find nearby restaurants, hotels or gas stations or to receive driving directions to a known address. Pioneer sells a car receiver that connects to an iPhone to take advantage of Siri® Eyes Free functionality including listening to text messages, calendar information, and reminders.\nA home automation system integrates electrical devices in a house with each other. The techniques employed in home automation include those in building automation as well as the control of domestic activities, such as home entertainment systems, houseplant and yard watering, pet feeding, changing the ambiance “scenes” for different events (such as dinners or parties), lighting control system, and the use of domestic robots. Devices may be connected through a home network to allow control by a personal computer, and may allow remote access from the internet. Through the integration of information technologies (e.g., PCs, phones) with the home environment, systems and appliances can communicate in an integrated manner which results in convenience, energy efficiency, and safety benefits.\nAutomated “homes of the future” have been staple exhibits for World's Fairs and popular backgrounds in science fiction. However, problems with complexity, competition between vendors, multiple incompatible standards, and the resulting expense have limited the penetration of home automation to homes of the wealthy or ambitious hobbyists.\nCurrent voice-controlled home automation systems typically involve wall-mounted control panels that provide a combination of touch and voice-controlled features. Such systems typically require a person to be standing next to the control panel, which substantially limits the ‘value add’ of using voice commands over using touch features of the control panel or a phone.\nTechnology improvements are needed to remove current limitations and inefficiencies of home (and business) automation systems."}
{"text": "In environments such as office buildings, shopping malls, airports, schools, sporting events, etc., personnel, which may include people generally, move between different locations to accomplish certain tasks or goals. For example, in an office building, in the event of a fire, personnel may move from their office location to an exit location with the goal of evacuating the office building. In an airport environment, personnel, which may include employees and travelers generally, may move from an airport entry to a check-in terminal, to security, and then to their specified gate for boarding and/or assisting with the boarding of an aircraft."}
{"text": "Heretofore a method of mapping the passivation layers of integrated circuits was disclosed by the applicant herein in U.S. Pat. No. 4,777,364, entitled \"Defect Detection and Thickness Mapping of the Passivation Layer(s) of Integrated Circuits,\" issued Oct. 11, 1988. In accordance with this patent, a penetration voltage method is used to determine the energy required by an electron beam to penetrate a passivation layer of an integrated circuit. As disclosed therein, an accelerating electron beam voltage is applied to an integrated circuit and is varied until the electrons have just enough energy to penetrate through the thickness of the passivation layer. Then, the acceleration voltage is increased to a predetermined amount (i.e., at least 3 KeV) so that the electron beam penetrates into the integrated circuit thin film and thereby interacts with the sublayer or film material to generate X-rays. The X-rays, in turn, are detected by an energy dispersive X-ray analyzer. The increased intensity-electron beam is x/y \"raster\" scanned over the area of the film and the generated X-ray intensity is stored at each pixel and displayed on a cathode ray tube.\nWith this method, however, greater accuracy is required in order to perform reliability tests on VHSI circuits and to map thickness profiles of the passivations layer in three dimensions. This lack of accuracy is due to a lack of resolution in the x-y plane which is a result of the averaging techniques employed by this method and it is due to a lack of specific correction for absorption effects in thin films. Accordingly, it is an objective of those who design VHSI circuits to obtain detailed thickness profiles of integrated circuits. The present invention addresses this objective."}
{"text": "The present invention relates to an analog to digital converter, and to a reference voltage source for a pipeline analog to digital converter.\nModern digital signal processing circuits are of central importance to recent advances in telecommunications, human/computer interface technology, image processing, and many other technologies. Analog to digital converters (ADC\"\"s) form an essential link in the signal processing pathway at the interface between the analog and digital domains. Advances in ADC technology have increased the speed, lowered the cost, and reduced the power requirements of analog to digital converters, and resulted in a proliferation of ADC applications.\nAmong existing ADC technologies are flash ADC, successive approximation ADC, Sigma-Delta ADC, and pipelined ADC. Flash ADC is performed by a highly parallel comparison of an input analog signal to each of a set of reference voltages. Flash ADC can provide very high speed and accuracy at the cost of high component count and high power consumption.\nSuccessive approximation ADC uses one or a few comparators, operated iteratively, to yield high accuracy conversion with far fewer components than flash conversion. Successive approximation ADC, however, operates at much slower conversion rates than flash ADC.\nSigma-Delta converters provide high accuracy conversion by oversampling, but at conversion rates that are also significantly slower than flash conversion.\nPipeline ADC provides analog to digital conversion that, while slower than flash conversion, is faster than most other ADC architectures. Pipeline ADC\"\"s introduce a latency (delay) between analog signal input and digital signal output. Conversion throughputs of pipeline ADC\"\"s, however, approach those of flash converters. Unlike flash converters, for which component counts increase exponentially with converter resolution, the component counts of pipeline ADC converters increase linearly with resolution. Consequently, pipeline ADC converters are relatively compact, inexpensive, and power efficient. Accordingly, pipeline ADC\"\"s are widely used in portable signal processing apparatus.\nPipeline ADC\"\"s require stable, low noise, reference voltages for optimum operation. Preferably, these reference voltages are available at low cost in terms of chip real estate and power consumption.\nFIG. 1 illustrates an exemplary pipeline ADC in block diagram form. The FIG. 1 circuit is shown as a single ended ADC. In common practice, however, many pipeline ADC\"\"s are implemented as fully differential circuits. Nevertheless, single ended representation has been chosen for FIG. 1 so as to reduce the complexity of the diagram, and enhance clarity of the disclosure. The exemplary converter FIG. 1 includes a 10-bit pipeline ADC such as might be integrated on a single substrate with a CMOS Active Pixel Sensor (APS) array.\nThe pipelined ADC 100 includes a sample-and-hold stage 102 followed by 9 conversion stages 104. Each conversion stage 104 includes a coarse ADC 106 for analog to digital conversion of a stage input signal received at a stage input 108. The coarse ADC 106 produces a 1.5 bit digital output signal at an output 110. A 1.5 bit output includes two output bits adapted to output only three possible states, rather than the four states available on a full 2 bit output. Each conversion stage 104 also includes a coarse digital to analog converter (DAC) 112 adapted to receive the 1.5 bit digital output signal of the coarse ADC 106 and produce a corresponding analog output voltage at an analog output 114. The digital output of the ADC conversion stage is also coupled to a digital correction circuit 118 having a plurality of digital inputs 120 each coupled to a respective one of the 9 conversion stages 104. Each conversion stage 104 further includes a subtracting node 122 with first 124 and second 126 analog inputs, and an analog output 128. Also included in the ADC stage 104 is a high precision gain element (amplifier) 130 with a gain of two.\nOperation of the above-described conversion stage 104 is as follows: an analog stage input signal is received at an input 107 of the coarse ADC 106 and at the first (positive) input 124 of the subtracting node 122. The coarse ADC 106 produces a 1.5 bit output representing one of three possible values. This 1.5 bit output is applied to the digital input 113 of the coarse DAC 112 which, responsively, produces an analog output signal with a magnitude equal to one of three possible output signal values. As further discussed below, these three output signal values are +VR/4, 0, and xe2x88x92VR/4 where VR is a reference voltage of particular magnitude. The output signal of the coarse DAC is applied to the second (negative) input 126 of the subtracting node 122. The subtracting node 122 produces an output equal to an arithmetic difference between the magnitude of the analog inputs at its first and second input terminals. This difference, referred to as a residual, is then applied to an input 131 of the high-precision gain stage 130. The precision gain stage 130 produces an amplified residual output signal at its output 134 having a magnitude equal to two times the magnitude of the residual signal. This amplified residual signal is passed on to the input 108 of the next successive ADC stage 104. Meanwhile, the digital output of the coarse ADC is received by the digital correction circuit 118 and logically combined with the digital outputs of the other 8 conversion stages 104 to produce a 10 bit digital output for the pipeline ADC at the output 140 of the digital correction circuit 118.\nFIG. 2 is a schematic diagram showing additional detail of the ADC conversion stage 104 described above with respect to FIG. 1. Note that as in FIG. 1, the FIG. 2 circuit is a simplified (single ended) representation of a circuit more commonly implemented as a fully differential stage. Accordingly, one sees an input terminal 108, a coarse ADC stage 106 including first 202 and second 204 comparators each having a respective first input 206 coupled to the input terminal 108 and a respective second input 208 coupled to a respective source 210, 212 of a respective reference voltage. The first 202 and second 204 comparators have respective first 214 and second 216 outputs coupled to respective first 218 and second 220 inputs of a digital latch circuit 224.\nThe digital latch circuit 224 includes a control input 226 and a 2 bit wide digital output 228. A coarse DAC 112 includes a multiplexer 240 with a 2-bit wide digital control input 242, first 246, second 248, and third 250 analog inputs and an analog output 252. The digital control input 242 of the DAC is coupled to the digital output 228 of the latch 224. As is well known, the analog output 252 of the multiplexer is switchingly coupled to, and assumes the electrical potential of, one of the analog inputs 246, 248, 250 depending on a signal received at the digital input 242.\nThe precision gain circuit 130 includes a high-gain differential amplifier 130 with a positive input 260, a negative input 262, and an output 264. The positive input 260 of the amplifier 130 is coupled to a source of constant potential (e.g. ground potential 300). The negative input 262 of the amplifier is coupled to a first plate 270 of a first capacitor 272, and a second plate of a second capacitor 276. The negative input 262 of the amplifier is also switchingly coupled to source of ground potential 300 by means of a switching device 280. The first capacitor 272 has a third plate 282 switchingly alternately coupled to the output 264 of the amplifier 130 and to the input terminal 108 of the ADC converter stage 104. The second capacitor 276 has a fourth plate 284 switchingly alternately coupled to the input terminal 108 of the ADC converter stage 104, and the analog output 252 of the multiplexer 240. The first 272 and second 276 capacitors have equal capacitance. Accordingly, the gain of the gain stage is 2 when the first capacitor 272 is switched into the feedback circuit 290.\nEach conversion stage 104 of the pipeline ADC 100 requires respective sources of five electrical potentials: ground 300 (common node voltage in a fully differential system), +VR applied at input 246, xe2x88x92VR applied at input 250, +VR/4 210, and xe2x88x92VR/4 212.\nFIG. 3 shows a conventional reference circuit 400 for generating the delta xe2x88x92Vref (=Vref_hixe2x88x92Vref-lo) differential reference voltage required by a fully differential pipeline ADC. The differential reference voltage deltaxe2x88x92Vref corresponds to the +VR and xe2x88x92VR reference voltages applied at the inputs 246 and 250 of the multiplexer 240 of the single-ended FIG. 2 circuit. The +Vref/4 and xe2x88x92Vref/4 signals required at the respective second inputs 208 of the FIG. 2 comparators 202, 204 are readily derived by a capacitive voltage dividing circuit, as known in the art. The corresponding reference voltages (delta xe2x88x92Vref/4) required by a fully differential pipeline ADC are achieved in the same manner.\nThe FIG. 3 circuit includes a fixed current source 404 coupled between a source of supply voltage 406 and one end 414 of a resistive ladder 408. The current source 404 is adapted to drive a fixed current through the resistive ladder 408. The resistive ladder includes a plurality of resistors 410 with a respective plurality of tap nodes 412 disposed therebetween. A second end 416 of the resistive ladder 408 is coupled to a source of ground potential 300. A first amplifier circuit 440 having a first (positive) 442 and a second (negative) 444 input and a first output 446 is provided. Also provided is a second amplifier circuit 450 with third (positive) 452 and a fourth (negative) 454 input and a second output 456. Both amplifier circuits 440, 450 are single ended.\nThe output 446 of the first amplifier circuit 440 is directly coupled back to the second negative input 444, yielding a gain of 1 for the first amplifier. The output 456 of the second amplifier circuit 450 is directly coupled back to the fourth negative input 454 yielding a gain of 1 for the second amplifier. The first 442 and third 452 inputs of the respective first 440 and second 450 amplifiers are coupled to respective output terminals 460, 462 of respective first 464 and second 466 switching devices. The first 464 and second 466 switching devices each has three inputs 480, each input 480 being coupled to a respective tap node 412 of the plurality of tap nodes.\nWhen electrical current is driven through the resistive ladder 408 by the current source 404, each tap node 412 assumes a particular electrical potential. When a particular tap node 412 is switchingly coupled to the respective positive input 442, 452 of the single ended amplifier 440, 450, the output 446, 456 of the amplifier assumes the voltage of the tap node. By an appropriate choice of tap nodes, a desired deltaxe2x88x92Vref can be established between the respective outputs 446, 456 of the first and second amplifiers. Because the first and second amplifiers are independent single-ended amplifiers, however, the voltage deltaxe2x88x92Vref between the output nodes 446, 456 is subject to common mode noise. Moreover, the current that flows through the current ladder dissipates substantial power. Reference circuit 400 is thus costly in terms of thermal budget and battery resources, particularly in the context of miniature equipment.\nAccordingly there is a need for a voltage reference circuit capable of supplying a stable and precise reference voltage deltaxe2x88x92Vref to an ADC circuit such as a fully differential pipeline ADC circuit.\nThe present invention applies a fully differential amplifier operating with negative feedback in the charge domain to source a stable reference voltage to, for example, a pipeline analog to digital converter (ADC). In one aspect, a reference voltage value is established by applying a voltage dropped across a single resistor to differential inputs of the fully differential amplifier through a sample and hold circuit.\nIn a further aspect of the invention the sample and hold circuit is a crowbar network adapted to transfer electric potential across a pair of matched capacitors in response to the closing of a crowbar switch. In a further aspect of the invention a fully differential amplifier is configured with capacitive feedback connections so as to exhibit unity gain.\nIn yet another aspect of the invention a control circuit is adapted to couple ancillary capacitors into the feedback circuit so as to controllably vary the gain exhibited by the differential amplifier.\nIn a still further aspect of the invention the differential output voltage output by the fully differential amplifier is divided with a further switched capacitor network to produce a divided reference voltage. In other aspect of the invention, the reference voltage and divided reference voltage are applied at inputs to an ADC such as a pipeline ADC used in a CMOS active pixel sensor array (APS).\nIn an additional aspect of the invention, the voltage dropped across the single resistor is established by the action of a programmable current source. The programmable current source is adapted to receive a control input corresponding to a required output current. The control input may be a numerical (digital) value or may be an analog signal, one or the other being implemented according to the requirements of a particular system. The control input may be a value received from an external controller, and may include a value based on a feedback signal taken from an output of the fully differential amplifier.\nThese and other advantages and features of the invention will be more clearly understood from the following detailed description of the invention which is provided in connection with the accompanying drawings."}
{"text": "Much research has been carried out in the field of human probiotics in the last decade (see review Huis in't Veld et al. (1994) Tibtech 12, 6–8). This research has been prompted by the rising interest by the public in their health and well-being. Many probiotic products are now available on the market and some of the beneficial effects derived from these products range from alleviation of lactose intolerance (Gilliland, S. E. (1990) FEMS Microbiol. Rev. 87, 175–188) to prevention of diarrheal diseases (Marteau, P. et al. (1993) FEMS Microbiol. Rev. 12, 207–220) and possible prevention of carcinogenesis (Adachi, S. (1992) In “The Lactic Acid Bacteria in Health and Disease”. (Wood, Ed.), 233–262, Elsevier, Barking). Controversy exists over many of these beneficial effects as no standardised procedures are available and contradictory results have been published with regard to the possible beneficial effects of cultured products containing ‘probiotic’ bacteria.\nPoor choice of strain has been cited as one of the contributing factors to the inconsistency and variability of results (Marteau, P. et al. (1993) supra) (Kim, H. S. (1988) Cult. Dairy Prod. J. 23, 6–9) and Fuller, R. ((1989) J. Appl. Bact. 66, 365–378) outlined criteria pertaining to the successful isolation of probiotic strains. The strains should be indigenous to the intended host species and also have the ability to (i) survive and grow within that host; (ii) exert a beneficial effect at the target site and (iii) be maintainable in the carrier food or system throughout product manufacture and storage.\nThere is a fast growing market for health-promoting products including probiotics. Many such products are now available (Jong, S. C. and Birmingham, J. M., (1993) ATCC Quart. Newslett. 13(1), 1–11). One of the more important components of these products is the microorganisms used. The most frequently utilised species include Bifidobacterium sp., Lactobacillus sp., and Propionibacterium sp. (O'Sullivan, M. G., et al. (1992) Trends in Food Sci. and Tech. 3(12), 309–314). There is a lack of substantiated evidence from controlled trials that the organisms currently used in such products are those which have beneficial effects on the gut flora (Tannock, G. W. (1983) In Human Intestinal Microflora in Health and Disease 517–5399 D. J. Hentges (ed.), New York, Academic Press). The source of the microorganism is critical to its survival and therefore its function in the human intestinal tract. Lee, Y-K and Salminen, S. ((1995) Trends Food Sci. Technol. 6, 241–245) stated that as a general requirement, a probiotic strain should be of human origin as some health-promoting effects may be species dependent. It is well known that the indigenous microflora is one of the major defense mechanisms that protects the human against colonisation by allochthonous invading bacteria (Tancrede, C. (1992) Eur. J. Clin. Microbiol. Infect. Dis. 11(11), 1012–1015) and it is also the human's best ally when supporting the immune system. Bacterial populations at different levels of the gastrointestinal tract constitute complex ecosystems depending on the physiology of the host and on interactions between bacteria.\nTen Brink et al. ((1994) Journal of Applied Bacteriology 77 140–148) isolated and screened a large number (˜1000) of Lactobacillus strains for the production of antimicrobial activity. Lactobacilli were isolated from various fermented foods and feeds (sauerkraut, cheese, sausage and silage), human dental plaque and faeces derived from different laboratory animals (rat, mouse, guinea pig and quail) and human volunteers. Only eight positive strains were found and two of these were studied, namely Lactobacillus salivarius M7 and Lactobacillus acidophilus M46. The former strain produces the broad spectrum bacteriocin salivaricin B which inhibits the growth of Listeria monocytogenes, Bacillus cereus, Brochothrix thermosphacta, Enterococcus faecalis and many lactobacilli. L. acidophilus M46 produces a bacteriocin acidocin B which combines the inhibition of Clostridium sporogenes with a very narrow activity spectrum within the genus Lactobacillus. However, these strains are not indigenous to the infected host species, which is one of the criteria which is required for a successful probiotic strain for human use.\nArihara, K. et al. ((1996) Letters in Applied Microbiology 22, 420–424) have isolated Salivacin 140 a bacteriocin from Lactobacillus salivarius subsp. salicinius T140. Strain T140 was isolated from the surface of Japanese pampas grass leaves grown close to an animal barn and thus the strain was likely to have derived from animal faeces.\nThere is a need for probiotic strains which meet the aforementioned criteria. Bacteriocin production by lactobacilli is thought to play an important role in the competitive exclusion of pathogens and other undesirable microorganisms of the intestinal tract of humans. Bacteriocins are broadly defined as proteinaceous compounds which exhibit a bactericidal effect against a wide range of microorganisms.\nDue to their diversity of species and habitats lactobacilli are the most bacteriocinogenic of the lactic acid bacteria. As many as forty bacteriocins produced by lactobacilli have now been isolated (Klaenhammer, T. R. (1993) FEMS Microbiol. Rev. 12, 39–86).\nBacteriocins have been isolated from human infant faeces. However, the bacteriocins were found to have narrow host ranges and were active only against other lactobacillus species (Toba, T. et al. (1991) Lett. Appl. Microbiol. 12, 228–231.).\nThere is a need for bacteriocins with a broad spectrum of activity."}
{"text": "This invention relates to the field of information networks, and more specifically to finding and distributing things to interested parties through a personalized network of trusted parties defined by a member.\nExisting methods of finding interested parties for things are inadequate. Most methods of finding interested parties are sent to loosely defined groups and do not allow the differentiation of parties things are available to. Selling things using such services as want ad Web sites and auction Web sites make things available to the entire general public that uses such services. Methods of targeting only certain parties require a great deal of manual labor to complete. To ensure selectivity in which person the thing is made available to, a person must personally ask close friends and family. This would be both time consuming and may not produce fruitful results because the limited size of the pool asked.\nExisting methods also do not allow their members to establish a priority between people or different groups of people. Finally, even if an interested party is found, there is difficulty in transferring the thing to the interested party because of varying methods of doing so (e.g., courier services, or mail).\nTherefore, what is needed is an electronic or other automated method of creating and storing personalized networks of trusted members on computerized networks."}
{"text": "1. Field of the Invention\nThis invention relates to filters and filtration material, and in particular, but not exclusively to tobacco smoke filters, and methods of producing same.\n2. Brief Description of Related Art\nNumerous methods of making filtration material have been proposed. One method previously proposed by the Applicant was disclosed in British Patent Specification No. 2 205 102A, wherein a particulate plastics material, a polysaccharide, optionally a binder, and water are fed to an extruder which is operated under such heat and pressure conditions that upon emergence from the extruder die, the extrudate assumes a cross-section greater than that of the exit orifice of the die. The expanded product may be fed to the garniture of a cigarette filter making machine or shredded and then gathered and formed into a cigarette filter in a garniture of a cigarette filter making machine. A disadvantage with the product obtained by following the teaching in this document was that, although the product was considerably expanded, it was not very suitable for standard filter making product because at low moisture contents the product could be friable and brittle. Thus, the use of such product on a filter tipping machine, where the filter element is rolled to interattach same to the tobacco rod, would be likely to result in the physical breakage thereof. Furthermore, it has been found that polypropylene, and other plastics materials having similar melting points, fuse with itself and/or the polysaccharide expansion medium and will not break down if left in the physical environment without leaving a fused amalgamation of plastics material."}
{"text": "This disclosure relates to a new and distinct variety (cultivar) of sweetpotato named ‘CX-1’.\nThis new variety of sweetpotato (Ipomoea batatas (L.) Lam.) was developed by selection of ‘Xushu 18’, and designated ‘CX-1.’ It was selected for large roots, very high dry matter content and higher yield. Field tests were conducted in Florida and South Carolina. Although there are a large number of sweetpotato varieties, there is a wide range of appearance, growth parameters, and uses, not a lot is known about the genetics of this crop species. Sweetpotato is the 7th largest food crop in the world and is a major food staple in many countries, mostly in the tropics.\nSweetpotatoes are underground storage roots, not tubers. The botanical description of a tuber is a short, thickened portion of an underground stem. The tuber has eyes composed of a ridge bearing a ‘scale-like’ leaf with tiny meristematic buds in the axial of this scale-like leaf. The sweetpotato is a true root, not an underground tuber that has buds. The anatomy of a storage root is the same as any root, with an additional ability to expand radially to store additional starch and nutrients. The sweetpotato crop is asexually propagated crop from ‘seed’ roots and plant cuttings.\nCurrently grown U.S. sweetpotato varieties were developed to be eaten as a highly nutritious and delectable vegetable. This variety type has been bred for specific traits for that use. They produce relatively uniform, attractively shaped storage roots, with dark orange flesh, a sweet, delicious flavor, and a moist texture (from 77% to 81% moisture). The typical size of these roots is from 6 to 16 ounces. The typical fresh weight yields of these varieties will range from 12 to 15 tons of marketable roots per acre.\nThe sweetpotato variety Xushu 18 is a variety that was publicly released in 1972 from China. It was bred by researchers in Jiangsu Province, China. Xushu 18 was a seedling from a cross from the variety Xindazi and breeding line 52-45. (Gitoner, C. 1996. Potato and Sweetpotato in China: systems, constraints, and potential. International Potato Center. Lima, Peru).\n‘Xushu 18’ has been a popular variety in China, a country that is a major producer of sweetpotatoes. Researchers in China report that during several decades, dozens of new varieties derived by mutation from named varieties have been registered. Researchers in Xuzhou, China report that a single plant resulting from a selection of ‘Xushu 18’ with superior characteristics was named ‘Xu 77-6’ and was selected to replace normal ‘Xushu 18’ (Daifu, M., Hongmin, L., DaPeng, L., and Yi, W. 2000. Sweetpotato varieties decline in China and the present practices. International Workshop on Sweetpotato Decline Study Sep. 8-9 2000. Kyushu National Agricultural Experiment Station. Miyakonojo, Japan.)\n‘Xushu 18’ was not developed as a table vegetable. It was developed for its production of starch and for animal feed. It is not sweet or delicious, but is white fleshed, very bland and very dry,—about 28% dry matter. Commercial yields in China have been reported from 6.88 ton up to 18.5 tons of fresh root weight per acre (Gitoner, 1996). Currently there are no sweetpotato varieties, with high dry matter content, commercially grown in the U.S. for industrial uses such as starch and ethanol.\nStarting with the ‘Xushu 18’, a distinctive, new cultivar was selected and developed for large roots, high dry matter and higher total yield. The original ‘Xushu 18’ cultivar was obtained from the Plant Genetic Resourses Unit, USDA, ARS, Griffin, Ga. In Spring 2008 ‘CX-1’, was first asexually propagated as cuttings from a single plant selection of ‘Xushu 18’. At the test farm in Colleton County, S.C., a single plant was selected for very high yield and propagated there for multiple location testing. After years of field trials, it was determined that a selected high-yielding single plant of Xushu 18, named ‘CX-1’, produced a higher yield of sweetpotatoes with a higher dry matter content compared to published data on ‘Xushu 18’ (Table 1).\nTABLE 1Fresh Fresh Weight inWeight inRoot DryResearcherLocationStatusTons/haTons/AMatter %Ma Daifu, China Virus 4017.829CIP(mean 19 Freelocations)uVirus 3415.228CheckGruneberg, PeruMean3013.625Crop SciTingo MariaHighest7031.5NRaLa MolinaLowest 5 2.2NRXie Yi Zhi, ThailandHighest2712.026ARCRyan-BohacUSAFL49.033‘CX-1’USASC37.533aNR is Not Reported;ha = hectors;A = acresWolfgang, G., Manrique, K., Zhang, D, and Hermann, M. 2005. Genotype X Environment interactions for a diverse set of sweetpotato clones evaluated across varying ecogeographic conditions in Peru. Crop Science. 45:2160-2171.Zhi, X. Y. Effect of Potassium on the yield of three sweet potato varieties. 1991. Asian Regional Vegetable Research Center. Proceedings of Training Workshop.Daifu, M., Hongmin, L., DaPeng, L., and Yi, W. 2000. Sweetpotato varieties decline in China and the prevent practices. International Workshop on Sweetpotato Decline Study Sep. 8-9, 2000. Kyushu National Agricultural Experiment Station. Miyakonojo, Japan.\nSweetpotato roots of ‘CX-1’ can be used to produce both starch and fuel ethanol. Estimates of at least 1500 gallons per acre of sweetpotatoes were based on field tonnage per acre, and on laboratory tests of the amount of ethanol produced per ton of dry matter from storage roots. Sweetpotato roots of ‘CX-1’ were also used to make chips and fries. The high dry-matter ‘CX-1’ roots are more suited to making chips and fries than those made from the sweet, watery-textured vegetable types of sweetpotato like the leading U.S. variety, ‘Beauregard’, (19-20% dry matter). The dry matter content of the ‘CX-1’ cultivar is at least 32%. In contrast, commercial vegetable type sweetpotatoes cultivars have a dry matter content that ranges from 19 to 23%. The texture of the ‘CX-1’ is dry and mealy in contrast to vegetable types that are moist and sweet with visible syrup in the flesh. The chips made from ‘CX-1’ were crisp, and light in color, similar to a commercial chip made from the white potato. (FIG. 5). Both the ‘CX-1’ chips and fries with high dry matter take up significantly less oil and were less bitter than these products made from the white potato.\nWhen properly cured and stored at 50° to 60° F., the ‘CX-1’ variety can store up to one year. Long-term storage is a genetically controlled trait and a requirement for any successful cultivar in the temperate zones of agricultural production like the U.S. This is in contrast to typical tropical dry fleshed varieties that will deteriorate very quickly in storage. The asexual reproductions run true to the original sweetpotato and to each other in all respects.\nSamples of the ‘CX-1’ variety will be deposited and maintained."}
{"text": "The current disclosure relates generally to foundations for mattresses to increase mattress support, and more particularly to foundations having novel slat construction to improve mattress support.\nWhile traditional foundations are sufficient in many situations, applicant believes that in some situations, improvement can be made. By way of example, in some situations, a sleeper's downward force on traditional slats may cause the mattress to dip in between the slats of the foundation (referred to herein generally as “mattress dip”). Applicant believes mattress dip may be caused by any of a number of problems, including at least one of mattress spring position relative to the slat, spring size, slat size, slat spacing, sleeper weight, sleeper habits, etc. Regardless of the cause, applicant believes mattress dip can result in decreased sleeping comfort and/or mattress life.\nAccordingly, applicant desires foundations that can reduce mattress dip by providing increased mattress support."}
{"text": "It has been heretofore conventional that the forming of a musical-tone signal for an electronic musical instrument is effected by forming a musical-tone signal waveform and passing the same through an envelope circuit to obtain a musical-tone signal having a specific envelope.\nThis envelope circuit requires a condenser and is defective in that the condenser may be rather large in size. As a result of the necessity of limiting the capacity for that reason, available envelopes are limited to those within a certain range and many envelopes which characterize various natural musical instruments or which have to be specially created cannot be formed as occasion demands."}
{"text": "The present invention pertains to a hardenable copolymer of tetrafluoroethylene (C.sub.2 F.sub.4), at least one other monomer selected from chlorotrifluoroethylene (C.sub.2 F.sub.3 Cl) or vinylidene fluoride (C.sub.2 H.sub.2 F.sub.2), an allylic polyol, and an allylic or acrylic monomer with a fluorocarbon chain. This copolymer is soluble in organic solvents and is particularly recommended for the fabrication of paints and varnishes whose principal properties include resistance to stains and soiling.\nThe fluorinated polymers are known for their good mechanical properties and their excellent resistance to chemical products and weather. However, their lack of solubility in conventional solvents prevents their use for certain applications such as resin for paints and varnishes where their properties are desired for the production of coatings with good resistance and easy maintenance.\nIn order to take advantage of the properties of the fluorinated polymers while avoiding their drawbacks, attempts have been make to make them soluble in the conventional organic solvents. In order to accomplish this, it is known to diminish the crystallinity of the fluorinated polymers by copolymerization of ethylenically unsaturated monomers, at least one of which is fluorinated. Because of its low crystallinity, such polymers generally have mediocre mechanical properties and, particularly, poor hardness. For this reason, it is desirable for certain applications, particularly when employing them in the fabrication of paints and varnishes, to preserve a sufficient degree of rigidity and to make them hardenable by incorporating functional groups in their structure.\nSuch hardenable fluorinated copolymers are described in French Patent No. 2,488,260. They have fluorine atoms at the alpha position of the principal chain which is indispensable for a product with good aging properties. The presence of fluorine atoms at the alpha position of the principal chain also brings to bear an effect of not retaining dust which is manifested by the fact that fluorinated paints soil less quickly then other paints. However, the presence of the fluorine at the alpha position of the principal chain is markedly insufficient to assure, in particular, protection against graffiti in the case of light-colored paints where a \"ghost\" image of the inscription will always subsist even after cleaning. In attempting to improve this stain-resistance characteristic, many solutions have been envisaged, principally involving the techniques described in the documents below.\nIn order to prevent coatings from becoming soiled quickly by dust, there is proposed in EP 186,186 a soil-resistant coating characterized by the combination of a fluorinated polymer that has reactive sites with specific fluorcarbon compounds of the type: EQU CF.sub.3 --(CF.sub.2).sub.n --CH.sub.2 --CH.sub.2 --OH EQU or EQU OH--CH.sub.2 --CH.sub.2 --(CF.sub.2).sub.m --CH.sub.2 --CH.sub.2 --OH\nThe addition of these products in small quantities generally does not have a stain-resistant effect. In larger amounts, they limit the cross-linking capability of the film and subsequently lead to the flotation of pigments in the paints and varnishes and thus have a negative impact on the regularity of the colors applied.\nIn EP 311,052 there is described a soil-resistant coating obtained by combination of a fluorinated polymer and a specific copolymer containing fluorine in the form of a polyfluorocarbon chain and molecular constituents from a hydrophillic monomer. The introduction of this specific copolymer does not have a favorable effect on the aging characteristics of the coating."}
{"text": "The desired characteristics for electrode materials which are to be used in redox cells to and other electrolytic processes include: low electrical resistivity, good mechanical properties and chemical stability and optionally low weight and low volume. The electrode material should not be attacked by reactants or products during cell operation. Thus there is a need for an electrode material which has good mechanical properties to withstand the hydraulic pressure and cell assembling operation without cracking and which is solution impermeable so that penetration and attack of the metal current collector (in end-electrodes) and loss in current efficiency (due to electrolytes crossing the bipolar electrodes) is low."}
{"text": "As Internet-connected devices become ever more ubiquitous, such devices are becoming physically smaller and otherwise resource-constrained. One term that has been used to describe this trend is the emerging “Internet of Things” (IoT) which utilizes resource-constrained devices to enable various types of internetworked functionality with respect to the objects or things in our everyday environment. In particular, such devices include sensors of various types including imaging sensors. For resource-constrained and particularly energy-constrained devices, compression of sensor data is desirable as a way to efficiently use resources such as wireless transmission resources. However, data compression is itself a computationally intensive task which can cause a significant and possibly excessive draw on device resources.\nIn addition to the challenges of resource-constrained devices, there are drawbacks associated with conventional computing architectures. In some computers, improvements in memory speed could not match the speed up gain by processor frequency due to increased focus on logic gate density, which resulted in a communication bottleneck. Hence, researchers introduced cache memories closer to the processor to help alleviate the bottleneck and improve the overall system speed. As further enhancement in frequency for a single core could not be achieved due to IO bandwidth limit, the focus moved to multi-core systems combined with memory hierarchy. Although such architecture supports parallel computing, the speed-up gain of the system heavily depends on the ability to parallelize the software code on that architecture. Also, putting accelerators closer to the main memory and sending the result of the executed task instead of the whole data has helped to mitigate the memory access bottleneck. However, the “von Neumann bottleneck” is still there since data needs to be sent back and forth between the memory and core blocks. Increase memory size and hierarchy improves system performance but has negative impact on energy and area. It also adds complexity especially for multi-core where memory coherency need to be dealt with. It is estimated that data migration from off-chip memory to the last level cache (LLC) and through the cache hierarchy, buses, and register file can account for up to 66% of the total energy in data-intensive applications such as those involving data compression.\nConventional attempts to solve the problems of implementing data compression for resource-constrained devices have shortcomings. For example, some conventional attempts are costly, inefficient (e.g., with respect to power consumption, resource utilization, and/or physical space utilization) and/or ineffective (e.g., with respect to particular performance criteria for particular applications). Embodiments of the invention address these and other problems as will be apparent from the description below."}
{"text": "To support progressive tumor growth beyond the size of 1-2 mm3, it is recognized that tumor cells require a functional stroma, a support structure consisting of fibroblast, smooth muscle cells, endothelial cells, extracellular matrix proteins, and soluble factors (Folkman, J., Semin Oncol, 2002. 29(6 Suppl 16), 15-8). Tumors induce the formation of stromal tissues through the secretion of soluble growth factors such as PDGF and transforming growth factor-beta (TGF-beta), which in turn stimulate the secretion of complimentary factors by host cells such as fibroblast growth factor (FGF), epidermal growth factor (EGF), and vascular endothelial growth factor (VEGF). These stimulatory factors induce the formation of new blood vessels, or angiogenesis, which brings oxygen and nutrients to the tumor and allows it to grow and provides a route for metastasis. It is believed some therapies directed at inhibiting stroma formation will inhibit the growth of epithelial tumors from a wide variety of histological types. (George, D. Semin Oncol, 2001. 28(5 Suppl 17), 27-33; Shaheen, R. M., et al., Cancer Res, 2001. 61(4), 1464-8; Shaheen, R. M., et al. Cancer Res, 1999. 59(21), 5412-6). However, because of the complex nature and the multiple growth factors involved in angiogenesis process and tumor progression, an agent targeting a single pathway may have limited efficacy. It is desirable to provide treatment against a number of key signaling pathways utilized by tumors to induce angiogenesis in the host stroma. These include PDGF, a potent stimulator of stroma formation (Ostman, A. and C. H. Heldin, Adv Cancer Res, 2001, 80, 1-38), FGF, a chemo-attractant and mitogen for fibroblasts and endothelial cells, and VEGF, a potent regulator of vascularization.\nA major regulator of angiogenesis and vasculogenesis in both embryonic development and some angiogenic-dependent diseases is vascular endothelial growth factor (VEGF; also called vascular permeability factor, VPF). VEGF represents a family of isoforms of mitogens existing in homodimeric forms due to alternative RNA splicing. The VEGF isoforms are reported to be highly specific for vascular endothelial cells (for reviews, see: Farrara et al. Endocr. Rev. 1992, 13, 18; Neufield et al. FASEB J. 1999, 13, 9).\nVEGF expression is reported to be induced by hypoxia (Shweiki et al. Nature 1992, 359, 843), as well as by a variety of cytokines and growth factors, such as interleukin-1, interleukin-6, epidermal growth factor and transforming growth factor. To date, VEGF and the VEGF family members have been reported to bind to one or more of three transmembrane receptor tyrosine kinases (Mustonen et al. J. Cell Biol., 1995, 129, 895), VEGF receptor-1 (also known as flt-1 (fms-like tyrosine kinase-1)), VEGFR-2 (also known as kinase insert domain containing receptor (KDR); the murine analogue of KDR is known as fetal liver kinase-1 (flk-1)), and VEGFR-3 (also known as flt-4). KDR and flt-1 have been shown to have different signal transduction properties (Waltenberger et al. J. Biol. Chem. 1994, 269, 26988); Park et al. Oncogene 1995, 10, 135). Thus, KDR undergoes strong ligand-dependant tyrosine phosphorylation in intact cells, whereas fit-1 displays a weak response. Thus, binding to KDR is believed to be a critical requirement for induction of the full spectrum of VEGF-mediated biological responses.\nIn vivo, VEGF plays a central role in vasculogenesis, and induces angiogenesis and permeabilization of blood vessels. Deregulated VEGF expression contributes to the development of a number of diseases that are characterized by abnormal angiogenesis and/or hyperpermeability processes. It is believed regulation of the VEGF-mediated signal transduction cascade by some agents can provide a useful mode for control of abnormal angiogenesis and/or hyperpermeability processes.\nThe vascular endothelial growth factors (VEGF, VEGF-C, VEGF-D) and their receptors (VEGFR2, VEGFR3) are not only key regulators of tumor angiogenesis, but also lymphangiogenesis. VEGF, VEGF-C and VEGF-D are expressed in most tumors, primarily during periods of tumor growth and, often at substantially increased levels. VEGF expression is stimulated by hypoxia, cytokines, oncogenes such as ras, or by inactivation of tumor suppressor genes (McMahon, G. Oncologist 2000, 5(Suppl. 1), 3-10; McDonald, N. Q.; Hendrickson, W. A. Cell 1993, 73, 421-424)\nThe biological activities of the VEGFs are mediated through binding to their receptors. VEGFR3 (also called Flt-4) is predominantly expressed on lymphatic endothelium in normal adult tissues. VEGFR3 function is needed for new lymphatic vessel formation, but not for maintenance of the pre-existing lymphatics. VEGFR3 is also upregulated on blood vessel endothelium in tumors. Recently VEGF-C and VEGF-D, ligands for VEGFR3, have been identified as regulators of lymphangiogenesis in mammals. Lymphangiogenesis induced by tumor-associated lymphangiogenic factors could promote the growth of new vessels into the tumor, providing tumor cells access to systemic circulation. Cells that invade the lymphatics could find their way into the bloodstream via the thoracic duct. Tumor expression studies have allowed a direct comparison of VEGF-C, VEGF-D and VEGFR3 expression with clinicopathological factors that relate directly to the ability of primary tumors to spread (e.g., lymph node involvement, lymphatic invasion, secondary metastases, and disease-free survival). In many instances, these studies demonstrate a statistical correlation between the expression of lymphangiogenic factors and the ability of a primary solid tumor to metastasize (Skobe, M. et al. Nature Med. 2001, 7(2), 192-198; Stacker, S. A. et al. Nature Med. 2001, 7(2), 186-191; Makinen, T. et al. Nature Med. 2001, 7(2), 199-205; Mandriota, S. J. et al. EMBO J. 2001, 20(4), 672-82; Karpanen, T. et al. Cancer Res. 2001, 61(5), 1786-90; Kubo, H. et al. Blood 2000, 96(2), 546-53).\nHypoxia appears to be an important stimulus for VEGF production in malignant cells. Activation of p38 MAP kinase is required for VEGF induction by tumor cells in response to hypoxia (Blaschke, F. et al. Biochem. Biophys. Res. Commun. 2002, 296, 890-896; Shemirani, B. et al. Oral Oncology 2002, 38, 251-257). In addition to its involvement in angiogenesis through regulation of VEGF secretion, p38 MAP kinase promotes malignant cell invasion, and migration of different tumor types through regulation of collagenase activity and urokinase plasminogen activator expression (Laferriere, J. et al. J. Biol. Chem. 2001, 276, 33762-33772; Westermarck, J. et al. Cancer Res. 2000, 60, 7156-7162; Huang, S. et al. J. Biol. Chem. 2000, 275, 12266-12272; Simon, C. et al. Exp. Cell Res. 2001, 271, 344-355). Moreover, VEGF activates the extracellular signal-regulated protein kinase (ERK) in human umbilical vein endothelial cells (HUVEC) (Yu, Y.; Sato, D. J. Cell Physiol 1999, 178, 235-246).\nPDGF is another key regulator of stromal formation which is secreted by many tumors in a paracrine fashion and is believed to promote the growth of fibroblasts, smooth muscle and endothelial cells, promoting stroma formation and angiogenesis. PDGF was originally identified as the v-sis oncogene product of the simian sarcoma virus (Heldin, C. H., et al., J Cell Sci Suppl, 1985, 3, 65-76). The growth factor is made up of two peptide chains, referred to as A or B chains which share 60% homology in their primary amino acid sequence. The chains are disulfide cross linked to form the 30 kDa mature protein composed of either AA, BB or AB homo- or heterodimmers. PDGF is found at high levels in platelets, and is expressed by endothelial cells and vascular smooth muscle cells. In addition, the production of PDGF is up regulated under low oxygen conditions such as those found in poorly vascularized tumor tissue (Kourembanas, S., et al., Kidney Int, 1997, 51(2), 438-43). PDGF binds with high affinity to the PDGF receptor, a 1106 amino acid 124 kDa transmembrane tyrosine kinase receptor (Heldin, C. H., A. Ostman, and L. Ronnstrand, Biochim Biophys Acta, 1998. 1378(1), 79-113). PDGFR is found as homo- or heterodimer chains which have 30% homology overall in their amino acid sequence and 64% homology between their kinase domains (Heldin, C. H., et al. Embo J, 1988, 7(5), 1387-93). PDGFR is a member of a family of tyrosine kinase receptors with split kinase domains that includes VEGFR2 (KDR), VEGFR3 (Flt4), c-Kit, and FLT3. The PDGF receptor is expressed primarily on fibroblast, smooth muscle cells, and pericytes and to a lesser extent on neurons, kidney mesangial, Leydig, and Schwann cells of the central nervous system. Upon binding to the receptor, PDGF induces receptor dimerization and undergoes auto- and trans-phosphorylation of tyrosine residues which increase the receptors' kinase activity and promotes the recruitment of downstream effectors through the activation of SH2 protein binding domains. A number of signaling molecules form complexes with activated PDGFR including PI-3-kinase, phospholipase C-gamma, src and GAP (GTPase activating protein for p21-ras) (Soskic, V., et al. Biochemistry, 1999, 38(6), 1757-64). Through the activation of PI-3-kinase, PDGF activates the Rho signaling pathway inducing cell motility and migration, and through the activation of GAP, induces mitogenesis through the activation of p21-ras and the MAPK signaling pathway.\nIn adults, it is believed the major function of PDGF is to facilitate and increase the rate of wound healing and to maintain blood vessel homeostasis (Baker, E. A. and D. J. Leaper, Wound Repair Regen, 2000. 8(5), 392-8; Yu, J., A. Moon, and H. R. Kim, Biochem Biophys Res Commun, 2001. 282(3), 697-700). PDGF is found at high concentrations in platelets and is a potent chemoattractant for fibroblast, smooth muscle cells, neutrophils and macrophages. In addition to its role in wound healing PDGF is known to help maintain vascular homeostasis. During the development of new blood vessels, PDGF recruits pericytes and smooth muscle cells that are needed for the structural integrity of the vessels. PDGF is thought to play a similar role during tumor neovascularization. As part of its role in angiogenesis PDGF controls interstitial fluid pressure, regulating the permeability of vessels through its regulation of the interaction between connective tissue cells and the extracellular matrix. Inhibiting PDGFR activity can lower interstitial pressure and facilitate the influx of cytotoxics into tumors improving the anti-tumor efficacy of these agents (Pietras, K., et al. Cancer Res, 2002. 62(19), 5476-84; Pietras, K., et al. Cancer Res, 2001. 61(7), 2929-34).\nPDGF can promote tumor growth through either the paracrine or autocrine stimulation of PDGFR receptors on stromal cells or tumor cells directly, or through the amplification of the receptor or activation of the receptor by recombination. Over expressed PDGF can transform human melanoma cells and keratinocytes (Forsberg, K., et al. Proc Natl Acad Sci USA., 1993. 90(2), 393-7; Skobe, M. and N. E. Fusenig, Proc Natl Acad Sci USA, 1998. 95(3), 1050-5), two cell types that do not express PDGF receptors, presumably by the direct effect of PDGF on stroma formation and induction of angiogenesis. This paracrine stimulation of tumor stroma is also observed in carcinomas of the colon, lung, breast, and prostate (Bhardwaj, B., et al. Clin Cancer Res, 1996, 2(4), 773-82; Nakanishi, K., et al. Mod Pathol, 1997, 10(4), 341-7; Sundberg, C., et al. Am J Pathol, 1997, 151(2), 479-92; Lindmark, G., et al. Lab Invest, 1993, 69(6), 682-9; Vignaud, J. M., et al, Cancer Res, 1994, 54(20), 5455-63) where the tumors express PDGF, but not the receptor. The autocrine stimulation of tumor cell growth, where a large faction of tumors analyzed express both the ligand PDGF and the receptor, has been reported in glioblastomas (Fleming, T. P., et al. Cancer Res, 1992, 52(16), 4550-3), soft tissue sarcomas (Wang, J., M. D. Coltrera, and A. M. Gown, Cancer Res, 1994, 54(2), 560-4) and cancers of the ovary (Henriksen, R., et al. Cancer Res, 1993, 53(19), 4550-4), prostate (Fudge, K., C. Y. Wang, and M. E. Stearns, Mod Pathol, 1994, 7(5), 549-54), pancreas (Funa, K., et al. Cancer Res, 1990, 50(3), 748-53) and lung (Antoniades, H. N., et al., Proc Natl Acad Sci USA, 1992, 89(9), 3942-6). Ligand independent activation of the receptor is found to a lesser extent but has been reported in chronic myelomonocytic leukemia (CMML) where the a chromosomal translocation event forms a fusion protein between the Ets-like transcription factor TEL and the PDGF receptor. In addition, activating mutations in PDGFR have been found in gastrointestinal stromal tumors in which c-Kit activation is not involved (Heinrich, M. C., et al., Science, 2003, 9, 9).\nCertain PDGFR inhibitors will interfere with tumor stromal development and are believed to inhibit tumor growth and metastasis.\nThe link between activity in tumor cell proliferation assays in vitro and anti-tumor activity in the clinical setting has been well established in the art. For example, the therapeutic utility of taxol (Silvestrini et al. Stem Cells 1993, 11(6), 528-35), taxotere (Bissery et al. Anti Cancer Drugs 1995, 6(3), 339), and topoisomerase inhibitors (Edelman et al. Cancer Chemother. \nCells protect their DNA by adopting a higher-order complex termed chromatin. Chromatin condensation is evident during mitosis and cell death induced by apoptosis while chromatin decondensation is necessary for replication, repair, recombination and transcription. Histones are among some of the DNA-binding proteins that are involved in the regulation of DNA condensation; and post-translational modifications of histone tails serve a critical role in the dynamic condensation/decondensation that occurs during the cell cycle. Phoshorylation of the tails of histone H3 is involved in both transcription and cell division (Prigent et al. J. Cell Science 2003, 116, 3677). A number of protein kinases have been reported to phosphorylate histone H3 and these kinases function both as signal transduction and mitotic kinases.\nPyrrolotriazine derivatives have been described as having kinase inhibitory activity in U.S. application Ser. No. 10/289,010, U.S. Pat. No. 6,670,357, WO 2001/19828, WO 2003/042172, WO 2004/009542, WO2004/009601, WO 2004/009784 and WO 2004/013145.\nIn one embodiment, the present invention provides a compound of formula (I)\nwhereinR1 is selected from the group consisting of aryl, benzyl, and heteroaryl,                wherein aryl and heteroaryl can be optionally substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl, wherein (C1-C4)alkyl can be substituted with 0, 1, 2 or 3 halogen, 0 or 1 heterocyclyl, or 0 or 1 (C1-C3)alkoxy, wherein                            (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,                                    (C1-C3)alkoxy, wherein (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,            halogen,            trifluoromethyl,            trifluoromethoxy,            (C3-C6)cycloalkyl,            phenyl optionally substituted with 1 or 2 halogen,                        \n                                     wherein X is CH2, O, S or NR1-1, and wherein R1-1 is hydrogen or (C1-C6)alkyl,            nitro,            cyano,            (C1-C3)alkylthio,            trifluoromethylthio,            (C1-C3)alkylcarbonyl,            (C1-C6)alkoxycarbonyl, and            phenoxy, wherein phenoxy can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of (C1-C6)alkyl, (C1-C6)alkoxy, trifluoromethoxy, and halogen,                        and        wherein benzyl can be substituted with 0, 1, 2 or 3 groups selected from halogen, (C1-C3)alkyl, and (C1-C3)alkoxy;R2 is selected from the group consisting of hydrogen, halogen, (C1-C4)alkyl and (C1-C4)alkoxy;R3 is selected from the group consisting of        carboxyl,        formyl,        (C1-C6)alkylcarbonyl optionally substituted with 0, 1, 2, or 3 groups selected from fluorine, chlorine, hydroxy, (C1-C6)alkoxy, and heterocycle,        (C3-C6)cycloalkylcarbonyl,        (C1-C6)alkoxycarbonyl optionally substituted with 0, 1, 2, or 3 groups selected from amino, and (C1-C6)alkoxycarbonyl,        aminocarbonyl,        (C1-C6)alkylaminocarbonyl, wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of (C3-C6)cycloalkyl, halogen, amino, (C1-C6)alkylamino, hydroxy, (C1-C6)alkoxy, (C1-C6)alkoxycarbonyl, (C1-C6)alkylthio, (C1-C6)alkoxycarbonylamino, and methylsulfonyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with or 0 or 1 heterocyclyl, wherein heterocyclyl can optionally be substituted with 0 or 1 (C1-C6)alkyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0 or 1 phenyl, wherein phenyl can optionally be substituted with 0 or 1 halogen, (C1-C6)alkyl, or (C1-C6)alkoxy,        heterocyclylcarbonyl optionally substituted with 0 or 1 amino, (C1-C6)alkylamino, cycloalkyl, or (C1-C6)alkyl, wherein (C1-C6)alkyl can optionally be substituted with 0 or 1 amino or (C1-C6)alkylamino,        (C1-C6)alkyl optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of                    a) hydroxyl,            b) amino,            c) (C1-C6)alkylamino, wherein (C1-C6)alkylamino can be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy, methylthio, and methylsulfonyl,            d) arylamino, wherein arylamino can be substituted with 0, 1 or 2 substituents independently selected from the group consisting of (C1-C6)alkyl, (C1-C6)alkoxy, and trifluoromethyl,            e) heterocyclyl, wherein heterocyclyl can be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1 or 2 hydroxy, methoxy or pyridyl,            f) imidazolyl,            g) pyridylamino,            h) (C1-C3)alkoxy optionally substituted by fluoro up to the perfluoro level, or by heterocycle, wherein heterocycle can optionally be substituted by 0 or 1 (C1-C6)alkyl,            i) (C1-C3)alkoxy(C2-C3)alkoxy, and            j) (C1-C6)alkoxycarbonyl,            k) (C3-C6)cycloalkyl,            l) cyano,                        (C1-C6)alkoxy optionally substituted with 1, 2 or 3 substituents independently selected from the group consisting of amino, (C1-C6)alkylamino, and heterocyclyl, wherein heterocyclyl can be substituted with 0, 1, 2 or 3 (C1-C6)alkyl,        (C3-C6)cycloalkylaminocarbonyl optionally substituted with (C1-C3)alkyl,        cyano,        heteroaryl, wherein heteroaryl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of                    a) (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1, 2, or 3 halogen, 0 or 1 heterocyclyl, 0 or 1 alkylamino, or 0 or 1 hydroxy or methoxy,            b) halogen,            c) amino,            d) alkylamino,            e) (C1-C6)alkoxycarbonyl, and            f) (C3-C6)cycloalkyl,                        heteroarylcarbonyl, which can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of (C1-C6)alkyl, (C3-C6)cycloalkyl and halogen,        heterocyclyl, wherein heterocyclyl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of (C1-C6)alkyl and (C1-C6)alkoxycarbonyl; andR4 is selected from the group consisting of hydrogen, (C1-C6)alkyl, (C1-C6)alkoxy and halogen;or a pharmaceutically acceptable salt thereof.        \nDepending on their structure, the compounds according to the invention can exist in stereoisomeric forms (enantiomers, diastereomers). The invention therefore relates to the enantiomers or diastereomers and to their respective mixtures. Such mixtures of enantiomers and/or diastereomers can be separated into stereoisomerically unitary constituents in a known manner.\nThe invention also relates to tautomers of the compounds, depending on the structure of the compounds.\nA salt for the purposes of the invention is a pharmaceutically acceptable salt of the compound according to the invention.\nPharmaceutically acceptable salts of the compounds (I) include acid addition salts of mineral acids, carboxylic acids and sulphonic acids, for example salts of hydrochloric acid, hydrobromic acid, sulphuric acid, phosphoric acid, methanesulphonic acid, ethanesulphonic acid, toluenesulphonic acid, benzenesulphonic acid, naphthalenedisulphonic acid, acetic acid, propionic acid, lactic acid, tartaric acid, malic acid, citric acid, fumaric acid, maleic acid and benzoic acid.\nPharmaceutically acceptable salts of the compounds (I) also include salts of customary bases, such as for example alkali metal salts (for example sodium and potassium salts, alkaline earth metal salts (for example calcium and magnesium salts) and ammonium salts derived from ammonia or organic amines having 1 to 16 carbon atoms, such as ethylamine, diethylamine, triethylamine, ethyldiisopropylamine, monoethanolamine, diethanolamine, triethanolamine, dicyclohexylamine, dimethylaminoethanol, procaine, dibenzylamine, N-methylmorpholine, dihydroabietylamine, arginine, lysine, ethylenediamine and methylpiperidine.\nSolvates for the purposes of the invention are those forms of the compounds that coordinate with solvent molecules to form a complex in the solid or liquid state. Hydrates are a specific form of solvates, where the coordination is with water.\nFor the purposes of the present invention, the substituents have the following meanings, unless otherwise specified:\nAlkyl represents a linear or branched alkyl radical having generally 1 to 6, 1 to 4 or 1 to 3 carbon atoms, representing illustratively methyl, ethyl, n-propyl, isopropyl, tert-butyl, n-pentyl and n-hexyl.\nAlkoxy in general represents a straight-chain or branched hydrocarbon radical having generally 1 to 6, 1 to 4 or 1 to 3 carbon atoms and bound via an oxygen atom. Non-limiting examples include methoxy, ethoxy, propoxy, isopropoxy, butoxy, isobutoxy, pentoxy, isopentoxy, hexoxy, isohexoxy. The terms “alkoxy” and “alkyloxy” can be used synonymously.\nAlkylamino represents an amino radical having one or two (independently selected) alkyl substituents, illustratively representing methylamino, ethylamino, n-propylamino, isopropylamino, tert-butylamino, n-pentylamino, n-hexylamino, N,N-dimethylamino, N,N-diethylamino, N-ethyl-N-methylamino, N-methyl-N-n-propylamino, N-isopropyl-N-n-propylamino, N-t-butyl-N-methylamino, N-ethyl-N-n-pentylamino and N-n-hexyl-N-methylamino.\nAlkylthio in general represents a straight-chain or branched hydrocarbon radical having 1 to 6 carbon atoms and bound via an sulfur atom. Non-limiting examples include methylthio and ethylthio.\nArylamino represents an amino radical having one or two (independently selected) aryl substituents, illustratively representing phenylamino.\nAminoalkyl represents an alkyl radical substituted with an amino group. Non-limiting examples include aminomethyl and aminoethyl.\nAminocarbonyl represents a free amide group.\nAlkylcarbonyl represents a carbonyl group having an alkyl substituent. Non-limiting examples include acetyl, propanoyl, butanoyl, 2-methylpropanoyl, and hexanoyl.\nAlkylaminocarbonyl represents an aminocarbonyl radical (free amide) having one or two (independently selected) alkyl substituents, illustratively representing methylaminocarbonyl, ethylaminocarbonyl, n-propylaminocarbonyl, isopropylamino-carbonyl, tert-butylaminocarbonyl, n-pentylaminocarbonyl, n-hexylaminocarbonyl, N,N-dimethylaminocarbonyl, N,N-diethylaminocarbonyl, N-ethyl-N-methylaminocarbonyl, N-methyl-N-n-propylaminocarbonyl, N-isopropyl-N-n-propylaminocarbonyl, N-t-butyl-N-methylaminocarbonyl, N-ethyl-N-n-pentylamino-carbonyl and N-n-hexyl-N-methylaminocarbonyl.\nAlkoxycarbonyl represents a carbonyl group having an alkoxy substituent. It illustratively represents methoxycarbonyl, ethoxycarbonyl, n-propoxycarbonyl, isopropoxycarbonyl, tert-butoxycarbonyl, n-pentoxycarbonyl and n-hexoxycarbonyl.\nCycloalkyl represents a monocyclic cycloalkyl radical having generally 3 to 8 or 5 to 7 carbon atoms, illustratively representing cyclopropyl, cyclobutyl, cyclopentyl, cyclohexyl and cycloheptyl.\nAryl represents a mono- to bicyclic carbocyclic radical, which is aromatic at least in one ring, having generally 6 to 10 carbon atoms, illustratively representing phenyl and naphthyl.\nHeteroaryl represents an mono- or bicyclic radical having generally 5 to 10 or 5 or 6 ring atoms and up to 5, in another embodiment up to 4 hetero atoms selected from the group consisting of nitrogen, oxygen and sulfur, which is aromatic at least in one ring. It can be attached via a ring carbon atom or a ring nitrogen atom. If it represents a bicycle, wherein one ring is aromatic and the other one is not, it can be attached at either ring. Illustrative examples are thienyl, furyl, pyrrolyl, pyrazolyl, thiazolyl, oxazolyl, imidazolyl, pyridyl, pyrimidyl, pyridazinyl, indolyl, indazolyl, benzofuryl, benzimidazolyl, benzothiophenyl, quinolinyl, isoquinolinyl, 1,3-benzodioxinyl, 1,4-benzodioxinyl, or benzodioxolyl.\nHeteroarylcarbonyl represents a heteroaryl residue bonded via a carbonyl carbon atom.\nHeterocyclyl represents a mono- or bicyclic, in another embodiment monocyclic, nonaromatic, i.e. saturated or partially unsaturated radical having generally 4 to 10 or 5 to 8 ring atoms and up to 3, in another embodiment up to 2 hetero atoms and/or hetero groups selected from the group consisting of nitrogen, oxygen and sulfur, CO, SO and SO2. If bicyclic, it can be a fused or spiro-connected bicycle. It can be attached via a ring carbon atom or a ring nitrogen atom. Illustrative examples are tetrahydrofuran-2-yl, pyrrolidin-2-yl, pyrrolidin-3-yl, pyrrolinyl, piperidinyl, morpholinyl, perhydroazepinyl.\nHeterocyclylcarbonyl represents a heterocyclyl residue bonded via a carbonyl carbon atom.\nMethylsulfonyl represents a —S(O)2CH3 residue.\nHalogen or halo represents a substituent selected from the group consisting of fluoro, chloro, bromo and iodo, in another embodiment fluoro and chloro.\n(C1-C3)alkoxy(C2-C3)alkoxy represents an 2 to 3 carbon alkoxy group substituted at the 2 or 3-position with a 1 to 3 carbon alkoxy group, illustratively representing 2-methoxyethoxy (CH3—O—CH2CH2—O—), 3-ethoxypropoxy (CH3CH2—O—CH2CH2CH2—O—), 2-methoxypropoxy (CH3—O(CH3)CHCH2—O—), 2-isopropoxyethoxy (CH3(CH3)CH—O—CH2CH2—O—) 2-methoxy-1-methylethoxy (CH3OCH2(CH3)CH—O—) or 3-propoxyethoxy, (CH3CH2CH2—O—CH2CH2—O—).\n(C3-C6)cycloalkylaminocarbonyl optionally substituted by (C1-C3)alkyl represents an aminocarbonyl radical (free amide) having one (independently selected) cycloalkyl substituent, which may be optionally substituted on the nitrogen atom by a (C1-C3)alkyl group and independently substituted on any available carbon atom by 1 or 2 (C1-C3)alkyl groups, illustratively representing N-cyclopropylaminocarbonyl, N-cyclopropyl-N-methylaminocarbonyl, N-(2-methylcyclopropyl)aminocarbonyl, N-cyclobutylaminocarbonyl, N-(2,2-dimethyl)cyclopropyl)aminocarbonyl, cyclopentylaminocarbonyl, N-(3-ethylcyclopentyl)aminocarbonyl, N-cyclohexylaminocarbonyl, N-cyclohexyl-N-ethylaminocarbonyl, N-(3-propylcyclohexyl)aminocarbonyl, and N-(4,4-dimethylcyclohexyl)aminocarbonyl.\nA * symbol next to a bond denotes the point of attachment in the molecule. Alternatively, a dotted line (---) denotes the bond via a radical is attached to the rest of the molecule.\nWhen prefixes such as (C1-C4) are used before substituents, they mean to indicate the respective number of carbon atoms, for example 1 to 4 in case of (C1-C4).\nWhen more than one substituent is selected from a group, selection can be independent of each other. Unless a maximum number of substituents is indicated, substitution can take place up to the maximum number of available substitution locations, e.g. in case of halogenation to the perhalo level.\nIn another embodiment, the present invention provides a compound of formula (I), wherein\nR1 is selected from the group consisting of phenyl and monocyclic heteroaryl having 5 or 6 ring atoms,\n                wherein phenyl and heteroaryl can be optionally substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl, wherein (C1-C4)alkyl can be substituted with 0, 1, 2 or 3 halogen, 0 or 1 pyrrolidinyl, 0 or 1 morpholinyl, or 0 or 1 (C1-C3)alkoxy wherein                            (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,                                    (C1-C3)alkoxy, wherein (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,            halogen,            trifluoromethyl,            trifluoromethoxy,            (C3-C6)cycloalkyl,            phenyl optionally substituted with 1 or 2 halogen,            trifluoromethylthio;R2 is selected from the group consisting of hydrogen, halogen, (C1-C4)alkyl and (C1-C4)alkoxy;R3 is selected from the group consisting of                        carboxyl,        formyl,        (C1-C6)alkylcarbonyl optionally substituted with 0, 1, 2, or 3 groups selected from fluorine, chlorine, hydroxy, (C1-C6)alkoxy, and monocyclic heterocycle having 5 or 6 ring atoms,        (C3-C6)cycloalkylcarbonyl,        (C1-C6)alkoxycarbonyl optionally substituted with 0, 1, 2, or 3 groups selected from amino, and (C1-C6)alkoxycarbonyl,        aminocarbonyl,        (C1-C6)alkylaminocarbonyl, wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of (C3-C6)cycloalkyl, halogen, amino, (C1-C6)alkylamino, hydroxy, (C1-C6)alkoxy, (C1-C6)alkoxycarbonyl, (C1-C6)alkylthio, (C1-C6)alkoxycarbonylamino, and methylsulfonyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0 or 1 hydroxyl or 0 or 1 monocyclic heterocyclyl having 5 or 6 ring atoms, wherein heterocyclyl can optionally be substituted with 0 or 1 (C1-C6)alkyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0 or 1 phenyl, wherein phenyl can optionally be substituted with 0 or 1 halogen or (C1-C6)alkyl,        monocyclic heterocyclylcarbonyl having 5 or 6 ring atoms, optionally substituted with 0 or 1 amino, (C1-C6)alkylamino, (C3-C6)cycloalkyl, or (C1-C6)alkyl, wherein (C1-C6)alkyl can optionally be substituted with 0 or 1 amino or (C1-C6)alkylamino,        (C1-C6)alkyl optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of                    a) hydroxyl,            b) amino,            c) (C1-C6)alkylamino, wherein (C1-C6)alkylamino can be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy, methylthio, and methylsulfonyl,            e) monocyclic heterocyclyl having 5 or 6 ring atoms, wherein heterocyclyl can be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1 or 2 hydroxy, methoxy or pyridyl,            f) imidazolyl,            h) (C1-C3)alkoxy optionally substituted by fluoro up to the perfluoro level, or by monocyclic heterocycle having 5 or 6 ring atoms, wherein heterocycle can optionally be substituted by 0 or 1 (C1-C6)alkyl,            i) (C1-C3)alkoxy(C2-C3)alkoxy, and            j) (C1-C6)alkoxycarbonyl,            k) (C3-C6)cycloalkyl,            l) cyano,                        (C3-C6)cycloalkylaminocarbonyl optionally substituted with (C1-C3)alkyl,        cyano,        heteroaryl, wherein heteroaryl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of                    a) (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1, 2, or 3 halogen, 0 or 1 monocyclic heterocyclyl having 5 or 6 ring atoms, 0 or 1 alkylamino, or 0 or 1 hydroxy or methoxy,            b) halogen,            e) (C1-C6)alkoxycarbonyl, and            f) (C3-C6)cycloalkyl,                        monocyclic heteroarylcarbonyl having 5 or 6 ring atoms,        monocyclic heterocyclyl having 5 or 6 ring atoms, wherein heterocyclyl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of (C1-C6)alkyl and (C1-C6)alkoxycarbonyl; andR4 is selected from the group consisting of hydrogen and halogen;or a pharmaceutically acceptable salt thereof.        \nIn another embodiment, the present invention provides a compound of formula (I), wherein\nR1 is selected from the group consisting of phenyl, pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, pyridinyl, and pyrimidinyl,\n                wherein phenyl, pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, pyridinyl, and pyrimidinyl can be optionally substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl, wherein (C1-C4)alkyl can be substituted with 0, 1, 2 or 3 halogen,            (C1-C3)alkoxy, wherein (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,            halogen,            trifluoromethyl,            trifluoromethoxy,            cyclopropyl,            phenyl optionally substituted with 1 or 2 halogen;R2 is selected from the group consisting of hydrogen, fluoro and chloro;R3 is selected from the group consisting of                        (C1-C6)alkylcarbonyl optionally substituted with 0, 1, 2, or 3 groups selected from fluorine, chlorine, hydroxy, (C1-C6)alkoxy, piperazinyl, morpholinyl, pyrrolidinyl, and piperidinyl,        cyclopropylcarbonyl,        aminocarbonyl,        (C1-C6)alkylaminocarbonyl, wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of (C3-C6)cycloalkyl, halogen, amino, (C1-C6)alkylamino, hydroxy, (C1-C6)alkoxy, (C1-C6)alkoxycarbonyl, (C1-C6)alkylthio, (C1-C6)alkoxycarbonylamino, and methylsulfonyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0 or 1 hydroxyl, piperazinyl, morpholinyl, pyrrolidinyl or piperidinyl, wherein piperazinyl, morpholinyl, pyrrolidinyl or piperidinyl can optionally be substituted with 0 or 1 (C1-C6)alkyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0 or 1 phenyl, wherein phenyl can optionally be substituted with 0 or 1 halogen or (C1-C6)alkyl,        heterocyclylcarbonyl selected from piperazinylcarbonyl, morpholinylcarbonyl, pyrrolidinylcarbonyl or piperidinylcarbonyl, optionally substituted with 0 or 1 amino, (C1-C6)alkylamino, (C3-C6)cycloalkyl, or (C1-C6)alkyl, wherein (C1-C6)alkyl can optionally be substituted with 0 or 1 amino or (C1-C6)alkylamino,        (C1-C6)alkyl optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of                    a) hydroxyl,            c) (C1-C6)alkylamino, wherein (C1-C6)alkylamino can be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of halogen, alkylamino, and methoxy,            e) piperazinyl, morpholinyl, pyrrolidinyl or piperidinyl, wherein piperazinyl, morpholinyl, pyrrolidinyl or piperidinyl can be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1 or 2 hydroxy or methoxy,            f) imidazolyl,            h) (C1-C3)alkoxy optionally substituted by fluoro up to the perfluoro level, or by monocyclic heterocycle having 5 or 6 ring atoms, wherein heterocycle can optionally be substituted by 0 or 1 (C1-C6)alkyl,            i) (C1-C3)alkoxy(C2-C3)alkoxy, and            j) (C1-C6)alkoxycarbonyl,            k) (C3-C6)cycloalkyl,            l) cyano,                        (C3-C6)cycloalkylaminocarbonyl optionally substituted with (C1-C3)alkyl,        cyano,        pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, isothiazolyl, pyridinyl, imidazolyl or pyrimidinyl, wherein pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, isothiazolyl, pyridinyl, imidazolyl or pyrimidinyl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of                    a) (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1, 2, or 3 halogen, 0 or 1 alkylamino, or 0 or 1 methoxy,            b) halogen, and            f) (C3-C6)cycloalkyl,                        pyrazolylcarbonyl, oxazolylcarbonyl, isoxazolylcarbonyl, thiazolylcarbonyl, pyridinylcarbonyl or pyrimidinylcarbonyl; andR4 is selected from the group consisting of hydrogen and fluoro;or a pharmaceutically acceptable salt thereof.        \nIn another embodiment, the present invention provides a compound of formula (I), wherein\nR1 is selected from the group consisting of phenyl and monocyclic heteroaryl,\n                wherein aryl and heteroaryl can be optionally substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl, wherein (C1-C4)alkyl can be substituted with 0, 1, 2 or 3 halogen, 0 or 1 pyrrolidinyl or morpholinyl, or 0 or 1 (C1-C3)alkoxy wherein                            (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,                                    (C1-C3)alkoxy, wherein (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,            halogen,            trifluoromethyl,            trifluoromethoxy,            (C3-C6)cycloalkyl,            (C1-C3)alkylthio, and            phenoxy, wherein phenoxy can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of (C1-C6)alkyl, (C1-C6)alkoxy, trifluoromethoxy, and halogen,R2 is selected from the group consisting of hydrogen, halogen, (C1-C4)alkyl and (C1-C4)alkoxy;R3 is selected from the group consisting of                        carboxyl,        (C1-C6)alkylcarbonyl optionally substituted with 0, 1, 2, or 3 groups selected from fluorine, chlorine, hydroxy, (C1-C6)alkoxy, and heterocycle,        (C3-C6)cycloalkylcarbonyl,        (C1-C6)alkoxycarbonyl,        aminocarbonyl,        (C1-C6)alkylaminocarbonyl, wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of (C3-C6)cycloalkyl, halogen, (C1-C6)alkylamino, hydroxy and (C1-C6)alkoxy, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with or 0 or 1 heterocyclyl, wherein heterocyclyl can optionally be substituted with 0 or 1 (C1-C6)alkyl, and wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0 or 1 phenyl, wherein phenyl can optionally be substituted with 0 or 1 halogen, (C1-C6)alkyl, or (C1-C6)alkoxy,        heterocyclylcarbonyl optionally substituted with 0 or 1 amino, (C1-C6)alkylamino, cycloalkyl, or (C1-C6)alkyl, wherein (C1-C6)alkyl can optionally be substituted with 0 or 1 amino or (C1-C6)alkylamino,        (C1-C6)alkyl optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of                    a) hydroxyl,            c) (C1-C6)alkylamino, wherein (C1-C6)alkylamino can be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy, methylthio, and methylsulfonyl,            e) heterocyclyl, wherein heterocyclyl can be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1 or 2 hydroxy, methoxy or pyridyl,            f) imidazolyl,            h) (C1-C3)alkoxy optionally substituted by fluoro up to the perfluoro level, or by heterocycle, wherein heterocycle can optionally be substituted by 0 or 1 (C1-C6)alkyl,            i) (C1-C3)alkoxy(C2-C3)alkoxy, and            j) (C1-C6)alkoxycarbonyl,            k) (C3-C6)cycloalkyl,                        (C3-C6)cycloalkylaminocarbonyl optionally substituted with (C1-C3)alkyl,        monocyclic heteroaryl having 5 or 6 ring atoms, wherein heteroaryl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of                    g) (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1, 2, or 3 halogen, 0 or 1 heterocyclyl, 0 or 1 alkylamino, or 0 or 1 hydroxy or methoxy,            h) halogen,            i) amino,            j) alkylamino,            k) (C1-C6)alkoxycarbonyl, and            l) (C3-C6)cycloalkyl,                        monocyclic heteroarylcarbonyl having 5 or 6 ring atoms,        monocyclic heterocyclylcarbonyl having 5 or 6 ring atoms, wherein heterocyclyl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of (C1-C6)alkyl and (C1-C6)alkoxycarbonyl; andR4 is hydrogen or fluoro;or a pharmaceutically acceptable salt thereof.        \nIn another embodiment, the present invention provides a compound of formula (I), wherein\nR1 is selected from the group consisting of phenyl, pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, pyridinyl, and pyrimidinyl,\n                wherein phenyl, pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, pyridinyl, and pyrimidinyl can be optionally substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl, wherein (C1-C4)alkyl can be substituted with 0, 1, 2 or 3 halogen,            (C1-C3)alkoxy,            halogen,            trifluoromethyl, and            phenoxy, wherein phenoxy can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of alkyl, (C1-C6)alkoxy, trifluoromethoxy, and halogen;R2 is hydrogen, fluoro, chloro, methyl, ethyl or methoxy;R3 is selected from the group consisting of                        (C1-C6)alkoxycarbonyl,        (C1-C6)alkylaminocarbonyl, wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of halogen, and        (C3-C6)cycloalkylaminocarbonyl optionally substituted with (C1-C3)alkyl,        piperazinylcarbonyl, morpholinylcarbonyl, pyrrolidinylcarbonyl or piperidinylcarbonyl; andR4 is hydrogen;or a pharmaceutically acceptable salt thereof.        \nIn another embodiment, the present invention provides a compound of formula (I), wherein\nR1 is selected from the group consisting of phenyl, pyrazolyl, thiazolyl, pyridinyl, and pyrimidinyl,\n                wherein phenyl, pyrazolyl, oxazolyl, isoxazolyl, thiazolyl, pyridinyl, and pyrimidinyl can be optionally substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    methyl, ethyl, propyl or butyl, wherein methyl, ethyl, propyl or butyl can be substituted with 0, 1, 2 or 3 fluoro or chloro,            fluoro or chloro,            trifluoromethyl, and            phenoxy, wherein phenoxy can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of methyl, ethyl, propyl or butyl, methoxy, ethoxy, propoxy, trifluoromethoxy, fluoro and chloro;R2 is hydrogen;R3 is selected from the group consisting of                        methoxycarbonyl, ethoxycarbonyl, propoxycarbonyl or butoxycarbonyl,        (C1-C4)alkylaminocarbonyl, wherein (C1-C4)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 fluoro, and        cyclopropylaminocarbonyl optionally substituted with methyl, ethyl or propyl,        pyrrolidinylcarbonyl or piperidinylcarbonyl; andR4 is hydrogen;or a pharmaceutically acceptable salt thereof.        \nIn another embodiment, the present invention provides a compound of formula (Ic)\nwhereinR1 is selected from the group consisting of aryl, benzyl, and heteroaryl,                wherein aryl and heteroaryl can be optionally substituted with 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl, wherein (C1-C4)alkyl can be substituted with 0, 1, 2 or 3 halogen, 0 or 1 pyrrolidinyl, 0 or 1 morpholinyl, or 0 or 1 (C1-C3)alkoxy wherein                            (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,                                    (C1-C3)alkoxy, wherein (C1-C3)alkoxy can be optionally substituted with (C1-C3)alkylamino,            halogen,            trifluoromethyl,            trifluoromethoxy,            (C3-C6)cycloalkyl,            phenyl optionally substituted with 1 or 2 halogen,                        \n                                     wherein X is CH2, O, S or NR1-1, and wherein R1-1 is hydrogen or (C1-C6)alkyl,            nitro,            cyano,            (C1-C3)alkylthio,            trifluoromethylthio,            (C1-C3)alkylcarbonyl,            (C1-C6)alkoxycarbonyl, and            phenoxy                        and        wherein benzyl can be substituted with 0, 1, 2 or 3 groups selected from halogen, (C1-C3)alkyl, and (C1-C3)alkoxy;R2 is selected from the group consisting of hydrogen, halogen, alkyl and alkoxy;R3 is selected from the group consisting of        carboxyl,        formyl,        (C1-C6)alkylcarbonyl optionally substituted with 1, 2, or 3 fluorine,        (C1-C6)alkoxycarbonyl,        aminocarbonyl,        (C1-C6)alkylaminocarbonyl, wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 1, 2 or 3 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy, and methylsulfonyl, and wherein (C1-C6)alkylaminocarbonyl can be substituted with 0 or 1 hydroxyl or 0 or 1 heterocyclyl,        heterocyclylcarbonyl,        (C1-C6)alkyl optionally substituted with 1, 2 or 3 substituents independently selected from the group consisting of                    a) hydroxyl,            b) amino,            c) (C1-C6)alkylamino, wherein (C1-C6)alkylamino can be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy, methylthio, and methylsulfonyl,            d) arylamino, wherein arylamino can be substituted with 0, 1 or 2 substituents independently selected from the group consisting of (C1-C6)alkyl, (C1-C6)alkoxy, and trifluoromethyl,            e) heterocyclyl, wherein heterocyclyl can be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1 or 2 methoxy or pyridyl,            f) imidazolyl,            g) pyridylamino,            h) (C1-C3)alkoxy optionally substituted by fluoro up to the perfluoro level,            i) (C1-C3)alkoxy(C2-C3)alkoxy, and            j) (C1-C6)alkoxycarbonyl,                        (C1-C6)alkoxy optionally substituted with 1, 2 or 3 substituents independently selected from the group consisting of amino, (C1-C6)alkylamino, and heterocyclyl, wherein heterocyclyl can be substituted with 0, 1, 2 or 3 (C1-C6)alkyl;        (C3-C6)cycloalkylaminocarbonyl optionally substituted with (C1-C3)alkyl,        cyano, and        heteroaryl wherein heteroaryl can be substituted with 0, 1, 2, or 3 groups independently selected from the group consisting of                    m) (C1-C6)alkyl, wherein (C1-C6)alkyl can be substituted with 0, 1, 2, or 3 halogen, 0 or 1 heterocycyl, 0 or 1 alkylamino, or 0 or 1 methoxy,            n) halogen,            o) amino, and            p) alkylamino;or a pharmaceutically acceptable salt thereof.                        \nIn another embodiment, the present invention provides a compound of formula (Ic):\nwhereinR1 is selected from the group consisting of aryl and heteroaryl,                wherein aryl and heteroaryl can be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of                    (C1-C4)alkyl,            (C1-C3)alkoxy,            halogen,            trifluoromethyl,            trifluoromethoxy,            (C3-C6)cycloalkyl,            phenyl optionally substituted with 0, 1 or 2 halogen, and                        \n                                     wherein X is CH2, O, S or NR1-1; and wherein R1-1 is hydrogen or (C1-C6)alkyl;R2 is selected from the group consisting of hydrogen, halogen, alkyl and alkoxy;R3 is selected from the group consisting of carboxyl, (C1-C6)alkoxycarbonyl, aminocarbonyl, (C1-C6)alkylaminocarbonyl, and heterocyclylcarbonyl,wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy and methylsulfonyl, orR3 is (C1-C6)alkyl optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of hydroxyl, amino, (C1-C6)alkylamino, arylamino, heterocyclyl, pyridyl, and pyridylamino,wherein (C1-C6)alkylamino can optionally be substituted with 0, 1, 2, 3 or 4 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy, methylthio, and methylsulfonyl, and wherein heterocyclyl can optionally be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can optionally be substituted with 1 or 2 methoxy or pyridyl,and wherein arylamino can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of (C1-C6)alkyl, (C1-C6)alkoxy, and trifluoromethyl, orR3 is (C1-C6)alkoxy optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of amino, (C1-C6)alkylamino, and heterocyclyl, wherein heterocyclyl can optionally be substituted with 0, 1, 2 or 3 (C1-C6)alkyl;or a pharmaceutically acceptable salt thereof.                        \nIn another embodiment, the present invention provides a compound of formula (I),\nwherein\nR1 is selected from the group consisting of phenyl and mono- or bicyclic heteroaryl containing 5, 6, 9 or 10 ring atoms and up to 2 hetero atoms selected from the group consisting of nitrogen, oxygen and sulfur,\n                wherein phenyl and heteroaryl can be substituted with 0, 1 or 2 substituents independently selected from the group consisting of (C1-C4)alkyl, C1-C3)alkoxy, halogen, trifluoromethyl, trifluoromethoxy, (C3-C6)cycloalkyl, phenyl optionally substituted with trifluoromethyl, and        \n                 wherein X is CH2, O, S or NR1-1;                    wherein R1-1 is hydrogen or (C1-C6)alkyl;R2 is selected from the group consisting of hydrogen, fluoro, chloro, methyl, and methoxy;R3 is selected from the group consisting of carboxyl, (C1-C6)alkoxycarbonyl, aminocarbonyl, and (C1-C6)alkylaminocarbonyl,wherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy and methylsulfonyl, orR3 is (C1-C6)alkyl substituted with amino, (C1-C6)alkylamino, pyridyl, or 5- to 6 membered heterocyclyl containing 1 or 2 heteroatoms selected from the group consisting of nitrogen, oxygen and sulfur,wherein (C1-C6)alkylamino can optionally be substituted with 0, 1 or 2 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy and methylsulfonyl, and wherein heterocyclyl can optionally be substituted with 0, 1 or 2 (C1-C6)alkyl, wherein (C1-C6)alkyl can optionally be substituted with 0 or 1 methoxy or pyridyl, orR3 is (C1-C6)alkoxy optionally substituted with 0, 1, 2 or 3 substituents independently selected from the group consisting of amino, (C1-C6)alkylamino, and a 5- to 6 membered heterocyclyl containing 1 or 2 heteroatoms selected from the group consisting of nitrogen, oxygen and sulfur, wherein heterocyclyl can optionally be substituted with 0, 1, 2 or 3 (C1-C6)alkyl;or a pharmaceutically acceptable salt thereof.                        \nIn another embodiment, the present invention provides a compound of formula (Ic),\nwhereinR1 is selected from the group consisting of                phenyl optionally substituted with 1-2 groups selected from halo, (C1-C4)alkyl, OCF3, CF3, and        \n                 wherein X is CH2, O, S or NR1-1, wherein R1-1 is hydrogen or (C1-C6)alkyl,        pyridyl optionally substituted with CF3,        pyrazolyl option substituted with 1-2 groups selected from (C1-C4)alkyl, (C3-C6)cycloalkyl, and phenyl optionally substituted with CF3,        isoxazolyl optionally substituted with (C1-C4)alkyl,        pyrimid-4-yl optionally substituted with (C1-C3)alkoxy        indazolyl, optionally substituted on N with (C1-C4)alkyl;R2 is hydrogen;R3 is selected from the group consisting of        CO2R3-1         CONR3-2R3-3         —(CH2)mNR3-4R3-5        \n                 wherein X is O or NH,        \n                 wherein X is O or NH,        \nm is 1, 2 or 3n is 1, 2, or 3;p is 1, 2 or 3;q is 2 or 3;r is 2 or 3s is 1, 2 or 3;t is 0, 1 or 2;u is 1, 2 or 3;R3-1 is H or (C1-C6)alkyl;R3-2 and R3-3 are independently selected from H and (C1-C6)alkyl;R3-4 and R3-5 are independently selected from H and (C1-C6)alkyl;R3-6 is CF3, (C1-C4)alkoxy or (C1-C4)alkyl;or a pharmaceutically acceptable salt thereof.\nIn another embodiment, the present invention provides a compound of the formula (Ia),\nwhereinR3 is selected from the group consisting of\n                 wherein X is O or NH,        \n                 wherein X is O or NH, and        aminocarbonyl or (C1-C6)alkylaminocarbonyl substituted as described above;n is 1, 2, or 3;R5 is independently selected from the group consisting of fluoro, chloro, methyl, ethyl, propyl, isopropyl, trifluoromethyl, trifluoromethoxy, and morpholino; andv is 0, 1 or 2;or a pharmaceutically acceptable salt thereof.        \nIn another embodiment, the present invention provides a compound of the formula (Ib),\nor a pharmaceutically acceptable salt thereof.\nIn another embodiment, the present invention provides a compound of the formula (Ic), wherein R1 represents\n\nIn another embodiment, the present invention provides a compound of the formula (Ic), wherein R3 represents ethoxycarbonyl.\nIn another embodiment, the present invention provides a compound of the formula (Ic), wherein R3 is (C1-C6)alkylaminocarbonyl,\nwherein (C1-C6)alkylaminocarbonyl can optionally be substituted with 1 or 2 substituents independently selected from the group consisting of halogen, amino, alkylamino, methoxy and methylsulfonyl, and wherein arylamino can optionally be substituted with 1 or 2 methoxy or trifluoromethyl.\nIn another embodiment, the present invention provides a compound as described as an Operational Example Examples section of the present application.\nIn another embodiment, the present invention relates to a compound capable of being metabolized or hydrolyzed to a compound of formula (I) under physiological conditions. This includes e.g. ester and amide derivatives (which can be hydrolyzed to the respective acids, alcohols and amines) as well as orthoesters and aminal esters (which can be hydrolyzed to the respective acids), acetals and hemiacetals (which can be hydrolyzed to the respective keto derivatives, e.g. the oxo group of the 4-oxopyrimidine ring moiety).\nIn another embodiment, the present invention provides a process for preparing a compound of formula (I), wherein a compound of formula (II)\nwherein R2, R3 and R4 have the meaning indicated above,is reacted with an isocyanate compound of formula (III)R1—NCO  (III)or with an carbamate of formula (VI)R1—NH—C(O)—OPh  (VI),wherein R1 has the meaning indicated above; ora compound of formula (IV)\nwherein R1, R2, R3 and R4 have the meaning indicated above,is reacted with an anine of formula (V)R1—NH2  (V),wherein R1 has the meaning indicated above.\nThe preparation of the compounds according to the invention can be illustrated by means of the following synthetic schemes. Unless specifically defined otherwise, the substituent placeholders such as R1 to R3 have the meaning indicated above.\nMethods for preparing pyrrolotriazines are also disclosed in published U.S. application Ser. No. 10/289,010 (Publication No. US 2003-0186982 A1), U.S. Pat. No. 6,670,357 (U.S. application Ser. No. 10/036,293), as well as WO 2003/042172, WO 2004/009542, WO2004/009601, WO 2004/009784 and WO 2004/013145, all of which are hereby incorporated by reference in their entirety.\nGeneral Methods of Preparation of Invention Compounds\nCompounds of the present invention of Formula (I) can be conveniently prepared from the corresponding amino compounds of Formula (II) by straightforward means as described in the Reaction Schemes below or by means well known to those skilled in the art. In these Reaction Schemes, unless otherwise specifically defined, the meanings of R1-R3, R3-2, R3-3, R3-4, R4 and r are identical to those described above.\nReaction Scheme 1 illustrates the general method of preparing Formula (I) compounds from the corresponding amino compounds of Formula (II) by standard methods of urea formation. In this scheme, a Formula (II) compound is allowed to react with either an isocyanate of Formula (III), or more preferably a carbamate of Formula (VI), generally in an inert solvent, to give the compound of Formula (I) directly. Alternatively, the amine of Formula (II) can be treated first with a chloroformate of Formula (VII), in an inert solvent, to provide an intermediate carbamate of Formula (IV). The Formula (IV) compound is then allowed to react with an amine of Formula (V), in an inert solvent, to provide the compound of Formula (I).\n\nA more specific example of the Reaction Scheme 1 method is illustrated in Reaction Scheme 2 below. In this scheme, the amine of Formula (II-1) [Formula (II), where R3 is CO2Et], is used as the starting material, reacting with an isocyanate ((III), carbamate (VI), or in a two step sequence using (VII) followed by (V), to give the compound of Formula (I-1) [Formula (I) where R3 is CO2Et]. The Formula (I-1) compound can then be used as a starting material to prepare other compounds of Formula (I) as shown, for example, in Reaction Scheme 3 below.\n\nReaction Scheme 3 illustrates the preparation of compounds of Formula (I) in which R3 is a variety of substituents, starting from the compound of Formula (I-1) in which R3 is CO2Et. For example, hydrolysis of (I-1), in either aqueous base or acid, provides the carboxylic acid compound of Formula (I-2). Coupling of this acid with an amine of Formula (R3-2)(R3-3)NH gives amide compounds of Formula (I-3).\nReduction of the Formula (I-1) ester with a reducing agent such as DIBAL, provides the alcohol of Formula (I-4). Oxidation of the alcohol by standard means such as the Dess-Martin periodinane gives the aldehyde of Formula (I-5). Conversion of the aldehyde to an amine compound of Formula (I-6) is accomplished by a reductive amination sequence. In this sequence, a primary amine of Formula R3-4—NH2 is added to the compound of Formula (I-5) in the presence of acetic acid, and the intermediate imine compound is not isolated but is selectively reduced with a reagent such as sodium triacetoxyborohydride to give the amine of Formula (I-6).\n\nThe alcohol of Formula (I-4) is further elaborated in Reaction Scheme 4, and the aldehyde of Formula (I-5) is further elaborated in Reaction Scheme 5, below.\nIn Reaction Scheme 4, the alcohol is converted to the homologous aldehyde by standard means, namely, conversion to a tosylate or mesylate with tosyl or mesyl chloride, respectively, and a base such as pyridine or Et3N, to give the intermediate of Formula (I-7). Reaction of (I-7) with a cyanide source, e.g., KCN or NaCN, in a polar solvent such as DMF gives the nitrile of Formula (I-8). Selective reduction with DIBAL with hydrolytic workup gives aldehyde of Formula (I-9). The Formula (I-9) aldehyde is converted to the compound of Formula (I-10) [(I) where R3 is (R3-4)NHCH2CH2—] by the reductive amination sequence as described for preparation of (I-6) in Reaction Scheme 3.\n\nThe aldehyde of Formula (I-5) serves as the starting material for the preparation of additional Formula (I) compounds as shown below in Reaction Scheme 5. A Wadsworth-Emmons type reaction of (I-5) with a phosphonic ester and strong base such as LiH gives the unsaturated ester of Formula (I-11); reduction of this ester to the saturated compound of Formula (I-12) is accomplished by hydrogenation using a platinum oxide catalyst in acetic acid. Reduction of the ester to the alcohol of Formula (I-13) is followed by oxidation to the compound of Formula (I-14), a 2-carbon homologue of the aldehyde of Formula (I-5). Reductive amination of (I-14), as previously described above in Reactions Schemes 3 and 4 provides the compound of Formula (I-15) [Formula (I) where R3 is (R3-4)NHCH2CH2CH2—].\nGeneral Methods of Preparation of Intermediates\nThe preparation of key intermediate (II-1) shown above as starting material for Reaction Scheme 2, is prepared as illustrated below in Reaction Scheme 6. A 4-nitrocinnamate of Formula (VIII) is allowed to react with the isocyanide reagent of Formula (IX) in the presence of a strong base such as lithium hexamethyldisilazide (LHMDS) in an aprotic solvent such as THF, to give the substituted pyrrole of Formula (X). Formylation of (X) under Vilsmeier conditions (e.g., DMF, POCl3) gives the 2-formylpyrrole of Formula (XI). The aldehyde (XI) is converted to the nitrile of Formula (XII) by reaction with hydroxylamine hydrochloride to form an intermediate oxime, which is dehydrated in situ to a nitrile of Formula (XII), using a reagent such as acetic anhydride. The nitrile of Formula (XII) is then N-aminated using a strong base such as NaH and an aminating reagent such as (Ph)2P(O)—O—NH2, to provide the N-amino nitrile of Formula (XIII). Reaction of (XIII) with formamide [HC(O)NH2] gives the pyrrolotriazine intermediate of Formula (XIV-1). Selective reduction of the nitro substituent of the phenyl ring is accomplished in the final step using a catalyst such as Raney-Nickel in THF, providing the intermediate (II-1). These same reduction conditions are used to convert addition compounds of general formula (XIV), prepared as shown in Reaction Schemes 8-11 below, to the corresponding formula (II) intermediates.\n\nThe cinnamates of Formula (VIII) are either commercially available or prepared as shown in Reaction Scheme 7. In this sequence, a substituted nitrotoluene of Formula (XV) is oxidized with a reagent such as potassium permanganate to give the corresponding acid of Formula (XVI); this acid is reduced to the alcohol of Formula (XVII) with a reducing agent such as borane and then oxidized to the aldehyde of Formula (XVIII) using a reagent such as the Dess-Martin periodinane. Wadsworth-Emmons type reaction of (XVIII) using (EtO)2P(O)CH2CO2Et and a strong base such as LiH gives the cinnamate of Formula (VIII).\n\nThe use of intermediate (XIV-1) for the preparation of the intermediate of Formula (II-2) is shown in Reaction Scheme 8 below. The Formula (XIV-1) compound is allowed to react with excess methyl Grignard reagent to give the tertiary alcohol of Formula (XIX). This compound is subjected to oxidative rearrangement and hydrolysis using hydrogen peroxide and a Lewis acid, such as BF3, to give the hydroxy compound of Formula (XX). Reaction of (XX) with a substituted alcohol of Formula (XXI) under Mitsunobu conditions, e.g., DEAD, TPP, gives the intermediate of Formula (XIV-2). Reduction of the nitro group in Formula (XIV-2), as described for preparation of Formula (II-1) in Reaction Scheme 6, gives the intermediate of Formula (II-2) [Formula (II) where R3 is R′R″N(CH2)rO— and R″ and R′ are as described in Reaction Scheme 8].\nThe compound of Formula (II-2) can be used to prepare the compound of Formula (I) [where R3 is a group of Formula R′R″N(CH2)rO—, and R″ and R′ are as described in Reaction Scheme 8] by the route outlined in Reaction Scheme 1.\n\nThe compound of Formula (XIV-1), prepared as shown above in Reaction Scheme 6, is used to prepare a variety of other intermediates of Formula (IV-3)-(XIV-7) as shown below in Reaction Scheme 9. For example, selective reduction of the ester group in Formula (XIV-1) using, for example, DIBAL in THF, gives the compound of Formula (XIV-3). Oxidation of the alcohol (XIV-3) to the aldehyde of Formula (XIV-4) is accomplished using standard conditions such as the Dess-Martin periodinane reagent in methylene chloride. The Formula (XIV-3) compound may also be converted to the corresponding chloride of Formula (XIV-5) with, for example, thionyl chloride. Reduction of both the nitro group and the chlorine in Formula (XIV-5) using Ra—Ni, provides the intermediate of Formula (II-3). Reaction of Formula (XIV-5) with an alcohol of Formula R3-9—OH and base such as sodium hydride gives the ether of Formula (XIV-6).\n\nAs shown below in Reaction Scheme 9, the compound of Formula (XIV-4) may also be used for the preparation of the nitrile of Formula (XIV-7), by a two step procedure: reaction with hydroxylamine hydrochloride and pyridine, followed by dehydration of the intermediate oxime using acetic anhydride. The Formula (XIV-4) aldehyde may also be converted to an isoxazole intermediate of Formula (XIV-8) by reaction with tosylmethylisocyanide (TosMIC) in the presence of a base such as potassium carbonate, in a protic solvent such as methanol.\n\nThe intermediate of Formula (XIV-1) may also be used for the preparation of other amides and heterocycles as shown in Reaction Scheme 11. Hydrolysis of (XIV-1) under standard conditions to the corresponding acid of Formula (XIV-9) is followed by conversion to the amide of Formula (XIV-11), either directly using an amine of formula R3-11-1—NH2, BOP, and a base such as TEA, or by prior conversion to the acid chloride of Formula (XIV-10), which is then allowed to react with the amine of formula R3-11-1—NH2. The acid chloride of Formula (XIV-10) is also used for the preparation of oxadiazoles of Formulae (XIV-12) and (XIV-13), by reaction with either 1) hydrazine and trimethyl orthoformate to give the oxadiazole of Formula (XIV-12), or 2) hydrazine, a carboxylic acid of formula R3-11-2—CO2H, and a dehydrating agent such as POCl3, to give the substituted oxadiazole Formula (XIV-13).\n\nThe intermediates of Formula (XIV-1)-(XIV-13) can then be reduced as desired to provide the corresponding amino compounds of Formula (II), and in turn the compounds of Formula (I).\n\n\nIt is also to be understood that starting materials are commercially available or readily prepared by standard methods well known in the art. Such methods include, but are not limited to the transformations listed herein.\nIf not mentioned otherwise, the reactions are usually carried out in inert organic solvents which do not change under the reaction conditions. These include ethers, such as diethyl ether, 1,4-dioxane or tetrahydrofuran, halogenated hydrocarbons, such as dichloromethane, trichloromethane, carbon tetrachloride, 1,2-dichloroethane, trichloroethane or tetrachloroethane, hydrocarbons, such as benzene, toluene, xylene, hexane, cyclohexane or mineral oil fractions, alcohols, such as methanol, ethanol or iso-propanol, nitromethane, dimethylformamide or acetonitrile. It is also possible to use mixtures of the solvents.\nThe reactions are generally carried out in a temperature range of from 0° C. to 150° C., preferably from 0° C. to 70° C. The reactions can be carried out under atmospheric, elevated or under reduced pressure (for example from 0.5 to 5 bar). In general, they are carried out under atmospheric pressure of air or inert gas, typically nitrogen.\nPro-drugs of this invention in general may be made by conventional methods well known in the art. For example, hydroxyl groups may be converted to esters by reacting the compounds with carboxylic acid chlorides or anhydrides under standard conditions. A hydroxyl group may also be converted to carbonates by reacting the compounds with chloroformates under standard conditions.\nSalts of the compounds identified herein can be obtained by isolating the compounds as hydrochloride salts, prepared by treatment of the free base with anhydrous HCl in a suitable solvent such as THF. Generally, a desired salt of a compound of this invention can be prepared in situ during the final isolation and purification of a compound by means well known in the art. Or, a desired salt can be prepared by separately reacting the purified compound in its free base form with a suitable organic or inorganic acid and isolating the salt thus formed. These methods are conventional and would be readily apparent to one skilled in the art.\nAdditionally, sensitive or reactive groups on the compound of this invention may need to be protected and deprotected during any of the above methods. Protecting groups in general may be added and removed by conventional methods well known in the art (see, for example, T. W. Greene and P. G. M. Wuts, Protective Groups in Organic Synthesis; Wiley: New York, (1999).\nIf used as active compounds, the compounds according to the invention are preferably isolated in more or less pure form, that is more or less free from residues from the synthetic procedure. The degree of purity can be determined by methods known to the chemist or pharmacist (see especially Remington's Pharmaceutical Sciences, 18th ed. 1990, Mack Publishing Group, Enolo). Preferably the compounds are greater than 99% pure (w/w), while purities of greater than 95%, 90% or 85% can be employed if necessary.\nThe compounds according to the invention exhibit an unforeseeable, useful pharmacological and pharmacokinetic activity spectrum. They are therefore suitable for use as medicaments for the treatment and/or prophylaxis of disorders in humans and animals.\nBecause of their antiproliferative properties, the compounds according to the invention are useful alone or in combination with other active components for treating and/or preventing mammalian hyper-proliferative disorders. Indications mediated by hyperproliferative disorders means diseases or conditions whose progression proceeds, at least in part, via proliferation.\nThe present invention also relates to a method of using the compounds or compositions described herein for the treatment or prevention of, or in the manufacture of a medicament for treating or preventing, mammalian hyper-proliferative disorders. This method comprises administering to a patient (or a mammal) in need thereof, including a human, an amount of a compound, a pharmaceutically acceptable salt or ester thereof, or a composition of this invention which is effective to treat or prevent the disorder.\nThe present invention also relates to a method for using the compounds of this invention as prophylactic or chemopreventive agents for prevention of the mammalian hyper-proliferative disorders described herein. This method comprises administering to a mammal in need thereof, including a human, an amount of a compound of this invention, or a pharmaceutically acceptable salt or ester thereof, which is effective to delay or diminish the onset of the disorder.\nHyper-proliferative disorders include but are not limited to solid tumors, such as cancers of the breast, respiratory tract, brain, reproductive organs, digestive tract, urinary tract, eye, liver, skin, head and neck, thyroid, parathyroid and their distant metastases. Those disorders also include lymphomas, sarcomas, and leukemias.\nThe present invention also relates to a method for using the compounds of this invention as prophylactic or chemopreventive agents for prevention of the mammalian hyper-proliferative disorders described herein. This method comprises administering to a mammal in need thereof, including a human, an amount of a compound of this invention, or a pharmaceutically acceptable salt or ester thereof, which is effective to delay or diminish the onset of the disorder.\nExamples of breast cancer include, but are not limited to invasive ductal carcinoma, invasive lobular carcinoma, ductal carcinoma in situ, and lobular carcinoma in situ.\nExamples of cancers of the respiratory tract include, but are not limited to small-cell and non-small-cell lung carcinoma, as well as bronchial adenoma and pleuropulmonary blastoma.\nExamples of brain cancers include, but are not limited to brain stem and hypothalamic glioma, cerebellar and cerebral astrocytoma, medulloblastoma, ependymoma, as well as neuroectodermal and pineal tumor.\nTumors of the male reproductive organs include, but are not limited to prostate and testicular cancer. Tumors of the female reproductive organs include, but are not limited to endometrial, cervical, ovarian, vaginal, and vulvar cancer, as well as sarcoma of the uterus.\nTumors of the digestive tract include, but are not limited to anal, colon, colorectal, esophageal, gallbladder, gastric, pancreatic, rectal, small-intestine, and salivary gland cancers.\nTumors of the urinary tract include, but are not limited to bladder, penile, kidney, renal pelvis, ureter, and urethral cancers.\nEye cancers include, but are not limited to intraocular melanoma and retinoblastoma.\nExamples of liver cancers include, but are not limited to hepatocellular carcinoma (liver cell carcinomas with or without fibrolamellar variant), cholangiocarcinoma (intrahepatic bile duct carcinoma), and mixed hepatocellular cholangiocarcinoma.\nSkin cancers include, but are not limited to squamous cell carcinoma, Kaposi's sarcoma, malignant melanoma, Merkel cell skin cancer, and non-melanoma skin cancer.\nHead-and-neck cancers include, but are not limited to laryngeal/hypopharyngeal/nasopharyngeal/oropharyngeal cancer, and lip and oral cavity cancer.\nLymphomas include, but are not limited to AIDS-related lymphoma, non-Hodgkin's lymphoma, cutaneous T-cell lymphoma, Hodgkin's disease, and lymphoma of the central nervous system.\nSarcomas include, but are not limited to sarcoma of the soft tissue, osteosarcoma, malignant fibrous histiocytoma, lymphosarcoma, and rhabdomyosarcoma.\nLeukemias include, but are not limited to acute myeloid leukemia, acute lymphoblastic leukemia, chronic lymphocytic leukemia, chronic myelogenous leukemia, and hairy cell leukemia.\nThese disorders have been well characterized in humans, and also exist with a similar etiology in other mammals which can also be treated by the administration of the compounds and/or pharmaceutical compositions of the present invention.\nThe assay described in this application is one of the methods by which compound activity relating to treatment of the disorders identified herein can be determined.\nIn another embodiment, the present invention provides a medicament containing at least one compound according to the invention. In another embodiment, the present invention provides a medicament containing at least one compound according to the invention together with one or more pharmacologically safe excipient or carrier substances, for example hydroxypropylcellulose, and also their use for the abovementioned purposes.\nThe active component can act systemically and/or locally. For this purpose, it can be applied in a suitable manner, for example orally, parenterally, pulmonally, nasally, sublingually, lingually, buccally, rectally, transdermally, conjunctivally, otically or as an implant.\nFor these application routes, the active component can be administered in suitable application forms. An overview of application forms is given in Remington's Pharmaceutical Sciences, 18th ed. 1990, Mack Publishing Group, Enolo.\nUseful oral application forms include application forms which release the active component rapidly and/or in modified form, such as for example tablets (non-coated and coated tablets, for example with an enteric coating), capsules, sugar-coated tablets, granules, pellets, powders, emulsions, suspensions, solutions and aerosols. Such sustained-release pharmaceutical compositions are described in Part 8, Chapter 91 of Remington's Pharmaceutical Sciences, 18th ed. 1990, Mack Publishing Group, Enolo.\nParenteral application can be carried out with avoidance of an absorption step (intravenously, intraarterially, intracardially, intraspinally or intralumbarly) or with inclusion of an absorption (intramuscularly, subcutaneously, intracutaneously, percutaneously or intraperitoneally). Useful parenteral application forms include injection and infusion preparations in the form of solutions, suspensions, emulsions, lyophilisates and sterile powders. Such parenteral pharmaceutical compositions are described in Part 8, Chapter 84 of Remington's Pharmaceutical Sciences, 18th ed. 1990, Mack Publishing Group, Enolo.\nIn one embodiment, the invention relates to intravenous (i.v.) application of the active compound, e.g. as bolus injection (that is as single dose, e.g. per syringe), infusion over a short period of time (e.g. for up to one hour) or infusion over a long period of time (e.g. for more than one hour). The application can also be done by intermittent dosing. The applied volume can vary dependent on the conditions and usually is 0.5 to 30, or 1 to 20 ml for bolus injection, 25 to 500, or 50 to 250 ml for infusion over a short period of time and 50 to 1000, or 100 to 500 ml for infusion over a long period of time.\nThe application forms have to be sterile and free of pyrogens. They can be based on aqueous solvents or mixtures of aqueous and organic solvents. Examples are ethanol, polyethyleneglycol (PEG) 300 or 400, aqueous solutions containing cyclodextrins or emulsifiers, such as lecithin, Pluronic F68®, Solutol HS15® or Cremophor®. Aqueous solutions are preferred.\nFor intravenous application the solutions are generally isotonic and euhydric, for example with a pH of 3 to 11, 6 to 8 or about 7.4.\nGlass or plastic containers can be employed as packaging for i.v.-solutions, e.g. rubber seal vials. They can contain liquid volumes of 1 to 1000, or 5 to 50 ml. The solution can directly be withdrawn from the vial to be applied to the patient. For this purpose, it can be advantageous to provide the active compound in solid form (e.g. as lyophilisate) and dissolve by adding the solvent to the vial directly before administration.\nSolutions for infusion can advantageously be packaged in containers made from glass or plastic, for example bottles or collapsible containers such as bags. They can contain liquid volumes of 1 to 1000, or 50 to 500 ml.\nForms suitable for other application routes include for example inhalatory pharmaceutical forms (including powder inhalers, nebulizers), nasal drops/solutions, sprays; tablets or capsules to be administered lingually, sublingually or buccally, suppositories, ear and eye preparations, vaginal capsules, aqueous suspensions (lotions, shake mixtures), lipophilic suspensions, ointments, creams, milk, pastes, dusting powders or implants.\nThe active components can be converted into the recited application forms in a manner known per se. This is carried out using inert non-toxic, pharmaceutically suitable excipients. These include inter alia carriers (for example microcrystalline cellulose), solvents (for example liquid polyethylene glycols), emulsifiers (for example sodium dodecyl sulphate), dispersing agents (for example polyvinylpyrrolidone), synthetic and natural biopolymers (for example albumin), stabilizers (for example antioxidants such as ascorbic acid), colorants (for example inorganic pigments such as iron oxides) or taste and/or odor corrigents.\nFor human use, in the case of oral administration, it is recommended to administer doses of from 0.001 to 100 mg/kg, or from 0.01 to 20 mg/kg. In the case of parenteral administration such as, for example, intravenously or via mucous membranes nasally, buccally or inhalationally, it is recommended to use doses of 0.001 to 0.60 mg/kg, in particular 0.01 to 30 mg/kg.\nIn spite of this, it can be necessary in certain circumstances to depart from the amounts mentioned, namely as a function of body weight, application route, individual behaviour towards the active component, manner of preparation and time or interval at which application takes place. It can for instance be sufficient in some cases to use less than the aforementioned minimum amount, while in other cases the upper limit mentioned will have to be exceeded. In the case of the application of larger amounts, it can be advisable to divide them into a plurality of individual doses spread through the day.\nThe percentages in the tests and examples which follows are, unless otherwise stated, by weight; parts are by weight. Solvent ratios, dilution ratios and concentrations reported for liquid/liquid solutions are each based on the volume."}
{"text": "The treatment of textile materials with various chemicals, dyestuffs, resins and the like has been long conducted by the use of aqueous baths. In such processes the fabric is essentially saturated by immersion in a water bath containing the treating chemical and eventually the water must be removed in order to continue the processing or to dry the fabric. Of the many procedures employed in the past for the treatment of fabrics, the most commonly employed is the pad-dry process in which the fabric is immersed and saturated with the aqueous treating solution, squeezed between rollers to a given wet pick-up and subsequently dried or dried and cured on a frame or heated drying roll before being taken up in a roll once again. The amount of water retained by the fabric is normally controlled by the pressure of the squeeze roll; in conventional methods a lower limit of about 50 to 70 percent water based on the weight of the fabric is still retained, depending upon the particular fabric used. This large amount of water requires a tremendous amount of energy in the form of heat to dry the fabric. It has been estimated that the amount of energy required to remove the water and dry the fabric is many times greater than the amount of energy that is needed in heating the cloth to carry out the desired chemical treating step, as for example, in the application and cure of a wash and wear finish on the fabric, or in the continuous dyeing of a fabric. In addition to the pad-dry process, in which the water is removed by squeezing between rollers, other procedures have recently been developed for more efficient removal of water. In one such procedure the saturated fabric is conveyed to a jet squeezer which employs a stream of compressed air jetting outward at the point of contact between the fabric and the nip rolls to substantially reduce the moisture content of the fabric. The use of this technique has resulted in a decrease of the water content in the fabric to about half of that normally remaining when using the squeeze roll technique discussed above. In another procedure vacuum extractor rolls are used. This process entails conveying the wet fabric as it exits from the treating bath over a perforated roll within which a vacuum is created whereby the moisture is extracted from the fabric. In some instances, roller coating methods can be used which continuously deliver aqueous treating composition to the fabric and the add-on is governed by the fabric speed and the rate of delivery of the treating composition by the coating roller. In this procedure the treating composition generally remains predominantly on or near the surface of the fabric, particularly when low add-ons are involved.\nWithin the past few years, several new approaches have been made to obtain uniform application of compositions to porous substrates. These recently developed procedures use foams in different form. However, the methods by which the foams had been applied to treat the fabric or yarn leave much to be desired. One such disclosure is to be found in U.S. Pat. No. 3,697,314 issued Oct. 10, 1972. In this patent there is shown a method for producing foam and then passing yarn through the foam so as to coat the exterior surface of the yarn with the foamed treating agent. It stresses that the yarn must pass through the foam agglomerate in order to assure a uniform distribution of the agent over the entire circumferential surface of the yarn as it passes through the foam. The reference shows no means by which the foam could be applied on only one surface of a fabric or material and still obtain uniform distribution or uniform penetration of the interior of the yarn or fabric. An earlier attempt to use foam for the treatment of textile materials is to be found in U.S. Pat. No. 1,948,568, issued Feb. 27, 1934. In this disclosure, a textile material is suspended in a closed container and foam is pumped into the container and forced through the textile material until the textile material is uniformally impregnated from all sides throughout the substrate structure and saturated with the textile treating agent in the form of a foam. In the batch process disclosed in this patent, the textile material is in a stationary or fixed position.\nThough disclosures do exist on the use of foam for the treatment of textile materials, essentially all of the industry still uses aqueous treating baths and processes in which the fabrics are generally immersed in the bath for the application of the treating material to the textile. As previously indicated, this entails the use of a large amount of energy to subsequently remove the water from the fabric. Further, none of the literature discloses foam compositions that have the adequate stability and penetration properties needed for treating a fabric without deposition of large quantities of water to it."}
{"text": "As an exemplary non-volatile semiconductor memory device, a type of non-volatile semiconductor memory device has been conventionally known that accumulates electric charges into a charge accumulation layer of a memory cell transistor by means of, for instance, a quantum tunneling effect in order to execute data writing (e.g., see PTL 1). Actually, as shown in FIG. 6, a non-volatile semiconductor memory device 100 of this type has a structure that higher order bit lines 101a and 101b and word lines 102a to 102h are disposed in an intersecting manner and a plurality of memory cell transistors 103 are disposed in a row and column matrix with respect to the higher order bit lines 101a and 101b and the word lines 102a to 102h. \nThe higher order bit line 101a is provided with a plurality of first semiconductor switches 104a and 104c, while a single lower order bit line 105a, 105c is connected to each first semiconductor switch 104a, 104c. Further, in this exemplary embodiment, the other higher order bit line 101b is also similarly provided with a plurality of first semiconductor switches 104b and 104d, while a single lower order bit line 105b, 105d is connected to each first semiconductor switch 104b, 104d. In such non-volatile semiconductor memory device 100, each lower order bit line 105a, 105b, 105c, 105d forms a memory block 106a, 106b, 106c, 106d, while each memory block 106a, 106b, 106c, 106d has the plural memory cell transistors 103.\nThe first semiconductor switches 104a, 104b, 104c and 104d are herein formed by N-MOS (Metal-Oxide-Semiconductor) transistors. Further, the first semiconductor switch 104a of the memory block 106a, for instance, is connected at the source thereof to the higher order bit line 101a while being connected at the drain thereof to the lower order bit line 105a, and is also connected at the gate thereof to a first selected gate line 108a shared with another memory block 106b aligned with the memory block 106a along a row direction. Thus, by this configuration, a predetermined gate voltage can be equally applied from the shared, single first selected gate line 108a to the two first semiconductor switches 104a and 104b that are mounted on the memory blocks 106a and 106b disposed in the upper part of FIG. 6.\nOn the other hand, similarly also in the two memory blocks 106c and 106d aligned along the row direction in the lower part of FIG. 6, a single first selected gate line 108b is connected to the two first semiconductor switches 104c and 104d, and a predetermined gate voltage is can be configured to be equally applied to the two first semiconductor switches 104c and 104d from the shared, first selected gate line 108b. \nIn addition to this, a higher order source line 110a is provided with a plurality of second semiconductor switches 111a and 111b, while a single lower order source line 112a, 112b is connected to each second semiconductor switch 111a, 111b. Further, the other higher order source line 110b is also similarly provided with a plurality of second semiconductor switches 111c and 111d, while a single lower order line 112c, 112d is connected to each second semiconductor switch 111c, 111d. \nFurther, the second semiconductor switches 111a, 111b, 111c and 111d are formed by NMOS transistors that the polarity thereof is the same as that of the first semiconductor switches 104a, 104b, 104c and 104d. \nHerein, the second semiconductor switch 111a of the memory block 106a, for instance, is connected at the source thereof to the higher order source line 110a while being connected at the drain thereof to the lower order source line 112a, and is also connected at the gate thereof to a second selected gate line 113a shared with another memory block 106b aligned along the row direction in the upper part. Thus, a predetermined gate voltage can be configured to be equally applied from the shared, single second selected gate line 113a to the two second semiconductor switches 111a and 111b that are mounted on the different memory blocks 106a and 106b disposed in the upper part.\nOn the other hand, similarly also in the two memory blocks 106c and 106d aligned along the row direction in the lower part, a single second selected gate line 113b is connected to the two second semiconductor switches 111c and 111d, and a predetermined gate voltage can be configured to be equally applied to the two second semiconductor switches 111c and 111d from the shared, second selected gate line 113b. \nFurther, each memory cell transistor 103 on the memory block 106a, for instance, is connected at one terminal thereof to the lower order bit line 105a while being connected at the other terminal thereof to the lower order source line 112a, and thus, the memory cell transistors 103 are disposed in parallel to each other between the lower order bit line 105a and the lower order source line 112a. The word lines 102a, 102b, 102c and 102d, shared by the memory block 106a and another memory block 106b aligned along the row direction, are connected to control gates of the memory cell transistors 103 of the memory block 106a. Due to this, a predetermined gate voltage can be configured to be equally applied, for instance, from the shared, single word line 102a to one of the memory cell transistors 103 of the memory block 106a in the upper part and one of the memory cell transistors 103 of another memory block 106b aligned with the memory block 106a along the row direction in the upper part.\nIncidentally, all the memory cell transistors 103 have the same structure that a channel region is disposed between one terminal and the other terminal, which are formed at a predetermined interval on a semiconductor substrate, and a charge accumulation layer, an interlayer insulation layer and a control gate are sequentially laminated through a tunnel insulation layer on the channel region of the semiconductor substrate. Such memory cell transistors 103 are of an N-channel type, and are capable of executing either data writing by injecting electric charges into the charge accumulation layer or data erasing by extracting the electric charges accumulated into the charge accumulation layer by means of voltage to be applied to the control gate and the region between the one terminal and the other terminal.\nThe non-volatile semiconductor memory device 100 thus structured can be configured to write data in a predetermined one of the memory cell transistors 103, read data from a predetermined one of the memory cell transistors 103 and erase data written in the memory cell transistors 103 by regulating voltages to be respectively applied to the higher order bit lines 101a and 101b, the higher order source lines 110a and 110b and the word lines 102a to 102h and by controlling on/off states of the first semiconductor switches 104a to 104d and the second semiconductor switches 111a to 111d. \nIn FIG. 6, the memory cell transistor 103 in the first row of the memory block 106a is set as a selected memory cell transistor 115 in which data is written, whereas all the remaining memory cell transistors 103 are set as non-selected memory cell transistors 116 in which data is not written.\nIt should be herein noted that for the sake of explanatory convenience, the memory block 106a on which the selected memory cell transistor 115 is disposed will be referred to as a selected block 117, whereas the memory blocks 106b, 106c and 106d on which only the non-selected memory cell transistors 116 are disposed will be referred to as non-selected blocks 118.\nActually in the non-volatile semiconductor memory device 100, for instance, when data is written in only the selected memory cell transistor 115 in the first row of the selected block 117, a high voltage of 12[V] is applied to a word line 120 (hereinafter referred to as a selected word line) that is the one connected to the selected memory cell transistor 115 among the plural word lines 102a to 102h, whereas a low voltage of 4[V] is applied to the word lines 121 (hereinafter referred to as non-selected word lines) that are the other remaining ones among the plural word lines 102a to 102h. \nFurther, at this time, in the non-volatile semiconductor memory device 100, a low voltage of 0[V] as a writing voltage can be applied to a higher order bit line 122 (herein referred to as a selected bit line) that is the one to which the selected memory cell transistor 115 is connected, whereas a high voltage of 8[V] as a writing prevention voltage can be applied to a higher order bit line 123 (herein referred to as a non-selected bit line) that is the one to which only the non-selected memory cell transistors 116 are connected. Moreover, in the non-volatile semiconductor memory device 100, a gate voltage of 10[V], which is higher than the voltage of the non-selected bit line 123, can be applied to the first semiconductor switches 104a and 104b from the first selected gate line 108a connected to the selected block 117, whereas a gate voltage of 0[V] can be applied to the second semiconductor switches 111a and 111b from the second selected gate line 113a. \nAccordingly, in the non-volatile semiconductor memory device 100, the first semiconductor switch 104b on the non-selected bit line 123 is switched on by means of the writing prevention voltage from the non-selected bit line 123 and the gate voltage from the first selected gate line 108a, and a writing prevention voltage of 8[V] can be applied to the non-selected memory cell transistor 116 on the non-selected bit line 123 that intersects with the selected word line 120. At this time, the second semiconductor switches 111a, 111b, 111c and 111d are switched off by applying a voltage of 0[V] thereto from the higher order source lines 110a and 110b and by applying a voltage of 0[V] thereto from the second selected gate lines 113a and 113b, and the lower order source lines 112a, 112b, 112c and 112d can be turned into a floating state.\nThus, in the non-selected memory cell transistor 116 at which the selected word line 120 and the non-selected bit line 123 intersect with each other, a voltage difference is reduced between the control gate and the semiconductor substrate, and as a result, electric charges cannot be injected into the charge accumulation layer without occurrence of a quantum tunneling effect.\nFurther, at this time, the first semiconductor switch 104a on the selected bit line 122 is switched on by means of the writing voltage from the selected bit line 122 and the gate voltage from the first selected gate line 108a, and a writing voltage of 0[V] can be applied to the selected memory cell transistor 115 on the selected bit line 122 that intersects with the selected word line 120. Accordingly, in the selected memory cell transistor 115 at which the selected word line 120 and the selected bit line 122 intersect with each other, a voltage difference is increased between the control gate and the semiconductor substrate by means of the writing gate voltage applied from the selected word line 120. As a result, a quantum tunneling effect occurs and electric charges can be injected into only the charge accumulation layer of the relevant selected memory cell transistor 115. Consequently, in the non-volatile semiconductor memory device 100, only the selected memory cell transistor 115 can be set to be in a data written state while electric charges are accumulated into the charge accumulation layer thereof."}
{"text": "The present invention relates to a cooking method and apparatus and in particular to such an apparatus specifically designed to cook shrimp.\nIn the past, a wide variety of cooking systems have been provided from simple pots to complex automated cooking ovens. Ovens having continuous conveyor belts have been provided in the past. These typically cook bakery products, and the like, and have electric or gas heating elements for baking cookies, bread, or the like, on conveyor belts passing through continuous ovens. This type of oven has a continuous flow of air at atmospheric pressure and is typically heated with electrical heat so that there is no build up of pressure as would be desirable in cooking seafood.\nShrimp are prepared utilizing a wide variety of prior art cooking methods and devices. In commercial cooking methods, however, it is generally customary to heat the shrimp with steam while applying pressure from an external source. See for example U.S. Pat. No. 3,501,317 to Veltman. In other methods such as the Hice U.S. Pat. No. 3,672,908, pressure from an external source is also applied prior to as well as during the heating of the shrimp.\nAn object of shrimp cooking processes is to avoid the weight loss during the heating step. As explained by Veltman, the cooking of shrimp in water, or an aqueous solutions, often causes a weight loss of between 30-50 percent resulting in a higher cost per pound of cooked shrimp.\nAnother object in any cooking process is to destroy or remove the bacteria present in the shrimp. Although acceptable for commercial purposes, many, if not most commercial shrimp cooking processes cannot remove bacteria below a 250,000 per gram plate count.\nThis invention is an improvement in my prior inventions set forth in U.S. Pat. No. 4,887,524 for a Shrimp Cooking Apparatus and my prior U.S. Pat. No. 4,966,072 for a Shrimp Cooking Apparatus. In those patents, a shrimp cooking apparatus used a frame having a motor driven endless belt conveyor with porous endless belts supported to the frame and a housing similar to the present invention. The steam is released directly beneath the belt and under the shrimp to heat the shrimp as they passed through the housing. In contrast, the present invention improves upon the prior systems by improving the housing which uses thin layers of stainless steel spaced by insulation to conserve energy in the system but which is reinforced with support brackets interconnected with a lifting mechanism formed in supporting upright channels and driven by central electric motors driving shafts to raise and lower the entire housing."}
{"text": "Bleeding in the gastrointestinal (“GI”) tract may be associated with various ulcers, lesions, cancers and the like. For example, peptic ulcers in the upper GI tract have been identified as a common cause of GI bleeding. If left untreated, GI bleeding may lead to anemia-like symptoms (e.g., fatigue, dizziness and chest pain), hepatic encephalopathy, hepatorenal syndrome, shock and death.\nSuccessful treatment of GI bleeding typically includes addressing the cause of the bleeding and/or haemostasis. For example, peptic ulcers may be associated with an infection of Helicobacter pylori and, therefore, may require treatment of the infection to reduce the risk of re-bleeding coupled with tissue coagulation to achieve haemostasis.\nHaemostasis may be achieved by invasive surgery or by various less invasive endoscopic techniques, such as laser treatment, bipolar electrocautery, heat probing, injections with sclerosing agents (e.g., epinephrine) or application of mechanical clips. While prior art endoscopic haemostasis techniques have presented some success, physicians continue to seek improved techniques for achieving haemostasis, while reducing damage to tissue adjacent to the treated tissue.\nAccordingly, there is a need for an improved apparatus and system for achieving haemostasis in the GI tract."}
{"text": "1. Field of the Invention\nThe present invention relates to a direct coupling type power amplifier for use with audio circuits and more particularly to such a direct coupling type power amplifier having a muting function.\n2. Description of the Prior Art\nIn a typical type of audio power amplifier circuit for amplifying audio signals, an inpact or popping sound unpleasant to the ears is sometimes generated from its load speaker immediately after the power source is turned on, due to a rapid change in bias voltages in a preamplifier such as an equalizer used in a first input amplifier stage of the power amplifier circuit. The impact sound is not only unpleasant to the ears, but also exposes the speaker to the danger of being damaged.\nThere has been provided a power amplifier circuit free from the generation of impact or popping sounds, in which a muting function is provided by breaking a signal transfer path through the actuation of an electromagnetic relay for a certain time interval after the turn-on of the power source. However, such a power amplifier circuit is expensive since the electromagnetic relay used therein is rather costly and moreover the switching sound of the relay is felt to be unpleasant."}
{"text": "1. Field of the Invention\nThe present invention relates to a power tool driven by a battery and, more particularly, to a power tool that can prevent an arc between the body of the power tool and the battery when the battery mounted is detached from the body during operation.\n2. Description of the Related Art\nAccording to a known power tool driven by a battery, the battery supplies driving current to a motor in order to drive a tool bit. For example, a known screwdriver includes a body having a motor housing and a hand grip, and a battery detachably coupled to the body. The motor housing includes the motor for driving the driver bit. The hand grip is connected to the motor housing. The battery is detachably coupled to the lower end portion of the hand grip. The battery is defined as one element of a power circuit to drive the motor by means of a connecting terminal. Thus, the battery supplies driving current to the motor via the power circuit.\nWithin the known power tool, because the battery is detachably coupled to the body, user of the power tool may unintentionally remove the battery from the body during the operation. Because the battery defines one element of the motor driving power circuit via a connecting terminal and the battery supplies driving current to the motor during its operation. Therefore, when the battery is removed from the body during the operation of the power tool, the power circuit may possibly be interrupted abruptly. When the power circuit for driving the motor is interrupted during the operation of the motor, arc may be generated at the contacts between the body of the power tool and the battery. Arc may adversely affect corrosion-resistance of the power tool."}
{"text": "1. Field of the Invention\nThe present invention relates to a varifocal lens for a TV camera which in particular may be mounted on an end scope, having a lens barrel, which is provided with guide indentations for radial pins may be moved toward each other for setting the focal length and together for focussing.\nBy way of illustration, TV cameras have increasingly been employed for taking endoscopic pictures. Such TV cameras must be as small and lightweight as possible in order not to hamper the physician during the operation. TV cameras with varifocal lenses, which would permit adjustment of the field of vision and adaption to various instruments have, for this reason, only been proposed now and again without having had any particular impact in practice.\nAn object of the present invention is to provide as small and light as possible a varifocal lens for a TV camera, which is particularly meant to be mounted on an endoscope.\nAn inventive solution to the object of the present invention and its further embodiments is described in the following.\nAn inventive step is that the lens barrel is guided in a support bearing, which is fixed in relation to the encasement of the TV camera, in such a manner that it can be moved in the direction of the optical axis of the lens and is secured to prevent rotary motion. Furthermore, the guide indentations are slits, which are open at the front and at the back and consist of a straight section abutting upon the front or rear end of the barrel and running parallel to the optical axis, a section running concentrically to the optical axis, and the hereto abutting actual cam for controlling the movement of the lens groups.\nAs a result the following advantages are yielded:\nThe lens mounting can be inserted into the lens barrel through the front or rear slits in such a manner that the lens barrel does not have to be constructed out of two half shells that consequently would have to be securely joined by means of a medium, thereby increasing the diameter. The support bearing, which prevents rotary motion while permitting moving the lens for focussing, does not increase the diameter of the varifocal lens.\nDue to the one-piece construction of the lens barrel, it suffices to arrange onto the lens barrel a second barrel in such a manner that it can be rotated, the inner surface of which is provided with indentations running parallel to the optical axis for the pins to engage. The focal length of the varifocal lens can easily be set by turning this second barrel.\nFor varying the focal length, the second barrel may, of course, be turned manually or by means of a setting lever.\nHowever, it is particularly advantageous if a tooth system is provided on the outer surface of the second barrel, which meshes together with a castellated wheel, as in this manner a simple knurled wheel on the outer side of the varifocal lens is all that is needed for setting the focal length permitting, as opposed to a setting ring that runs concentrically to the optical axis, easy one-handed operation for setting both the focal length as well as focussing. The knurled wheel for focussing may, by way of illustration, engages the eccentric, elliptical indentations in the second barrel by means of pins.\nThe stability of the slit lens barrel and a steadiness of the setting are substantially increased by means of the further embodiment wherein the lens mountings, on the one hand, give the lens barrel stability and, on the other hand, as a result of this embodiment they cannot tip.\nA further embodiment described in claim 4 hereto, which is provided with two lugs that run parallel to the optical axis, makes very compact assembly possible as the support bearing can be provided at the point of separation between the lens and the tube for the TV camera, where it takes up practically no additional room.\nYet a further embodiment, which is provided with knurled wheels for focussing as well as for setting the focal length, is ergonomically advantageous, as it makes one-hand operation possible without changing the lever position.\nA water protection rim also contributes to solving the object to of the present invention as no measures that would increase volume or weight, like those required according to the state of the art, are needed when flanging the varifocal lens to an endoscope in order to prevent water from running onto the front cover."}
{"text": "1. Field of the Invention\nThe present invention relates to an active matrix display apparatus having pixels each including a light-emitting device, and to a method for making the active matrix display apparatus.\n2. Description of the Related Art\nIn recent years, increasing efforts have been made to develop flat self-emission display apparatuses in which organic electroluminescent (EL) devices are used as light-emitting devices. An organic EL device is a device using a phenomenon in which an organic thin film emits light when an electric field is applied thereto. The organic EL device, which is driven by the application of a voltage of 10 V or less, is a low power consumption device. At the same time, since the organic EL device, which is a self-emission device capable of emitting light by itself, requires no illuminating unit, it is easy to produce an organic EL device that is thin and lightweight. Additionally, since the response speed of the organic EL device is as very high as several microseconds (μs), it is possible to prevent an afterimage from appearing when a moving image is displayed.\nOf flat self-emission display apparatuses with pixels each including an organic EL device, active matrix display apparatuses in which thin-film transistors (TFTs) are formed as drive devices in each pixel in an integrated manner have been particularly actively developed. For example, flat self-emission display apparatuses of active matrix type are described in the following documents: Japanese Unexamined Patent Application Publications Nos. 2003-255856, 2003-271095, 2004-133240, 2004-029791, 2004-093682, and 2005-166687.\nAn active matrix display apparatus of related art includes a substrate having wiring lines including signal lines arranged in columns, scanning lines arranged in rows, and predetermined power supply lines; and a matrix of pixels, each pixel being disposed at an intersection between a signal line and a scanning line. The wiring lines are formed by patterning a conductor film. Each pixel includes active devices (for example, TFTs) and a light-emitting device (for example, an organic EL device) which are connected to the wiring lines. The pixel operates in response to a control signal supplied from a scanning line. According to a video signal supplied from a signal line, the pixel causes a drive current supplied from a power supply line to flow through the light-emitting device."}
{"text": "Field of the Invention\nThis invention generally relates to a hydraulic hose fitting. More specifically, the present invention relates to a hydraulic hose fitting that is used with a hydraulic device (e.g., a hydraulic brake operating device, a hydraulic brake caliper, etc.).\nBackground Information\nIn recent years, some bicycles have been provided with hydraulic braking systems. The hydraulic braking system for a bicycle typically has a hydraulic brake operating device that is fluidly coupled to a hydraulic brake caliper by a hydraulic hose. The hydraulic brake caliper is hydraulically actuated by hydraulic fluid flowing through the hydraulic brake hose in response to operation of the hydraulic brake operating device. In particular, operation of the hydraulic brake operating device forces hydraulic fluid through the hydraulic hose to the brake caliper. The hydraulic fluid then moves one of more pistons to cause brake pads of the hydraulic brake caliper to squeeze a rotor that is attached to a hub of a bicycle wheel. Often, a hydraulic hose fitting is used at each end of the hydraulic hose to connect the hydraulic hose to the hydraulic brake operating device and the hydraulic brake caliper."}
{"text": "Services such as Google Maps (maps.google.com) allow users to conduct searches that are limited to particular geographic areas. For example, while viewing a street or satellite map of New York City on a web browser such as Google Chrome, a user may search for the word “pizza”. In response, the servers cause the user's computer to display icons indicating the positions of listings that match the query, such as restaurants that have “pizza” in their name.\nThe results of the query may be sent in a way that permits a great deal of interactivity. For example, the map is sent as tiles of images and descriptions of the results and their positions are sent as text (such as data formatted as JSON (Javascript Object Notation) or HTML). This allows the user's computer to display icons on the map and a text listing of the results next to the map. Using Javascript provided by the server, the client computer can determine whether the user has clicked an area of the map occupied by an interactive icon and perform a number of actions in response, such as displaying a balloon containing even more information about the listing.\nGoogle Maps also permits layers. For example, it is capable of showing icons for entries on Wikipedia.org that are associated with latitude/longitude coordinates matching the geographic area currently being viewed (e.g., http://en.wikipedia.org/wiki/Wikipedia:WikiProject_Geographica 1_coordinates). The icons of entries are sent as an image to be overlaid on the map, along with information relating to the entries. The information is sent in JSON with information identifying the entry's position on the map and an identification of the entry. The position information may be used by the user's computer to determine whether the user actuated the icon (such as by clicking it or hovering a mouse cursor over it) and the identification information may be used to obtain more information from a server if the icon is actuated. YouTube results may also be shown as a layer. These results are not filtered based on the query but, rather, show all entries in the area covered by the map.\nGoogle Maps also features a “Places of Interest” layer that allows users to display icons associated with all businesses falling within certain pre-selected categories sent by the server, such as “Italian” and “Pizza” restaurants, banks, gas stations, and the like. Although the user may collectively display search results from the pre-selected categories and from searches based on custom search terms that are not pre-selected, the results are separate, i.e., the preselected category icons do not change based on the user-defined term and vice versa. For example, the user may check a box next to “pizza” and type in the term “vegetarian”, in which case Google Maps combines and displays the search results together; however, the results associated with “pizza” are not selected based on the “vegetarian” query nor are the results associated with the “vegetarian” query selected based on the selection of “pizza” from the list. Sites such as www.tavernoxoros.gr (showing many locations on a map in response to a user's selection of pre-selected categories) and www.smalltownbrooklyn.com (which may be capable of showing all businesses on a map) are similarly able to show all results associated with pre-selected categories."}
{"text": "The present invention relates to structures and methods of placing carbon nanotubes (CNTs) on a substrate. More particularly, the present invention is generally directed to structures and methods that include selective electrostatic placement based on a dipole-to-dipole interaction of electron-rich carbon nanotubes onto an electron-deficient pre-patterned surface.\nCNTs can be semiconducting and therefore are of interest for electronic applications. Semiconducting CNTs can replace or complement traditional semiconductors in both high-performance and low-cost thin film transistor devices. To make any electronic device, one needs to precisely place the nanostructure over large areas. Various approaches have been developed for placing CNTs on a substrate."}
{"text": "The present invention relates in general to an apparatus for lifting and carrying large, heavy loads and, in particular, to an apparatus for moving objects up to and around on the roof of a building.\nMany commercial buildings are constructed with a relatively flat roof. Such roofs are typically covered with strips of roofing material which are delivered to the building site in large, heavy rolls. These rolls must be lifted to the roof of the building and then typically moved by hand by the workers installing the roofing material. Such an operation is difficult and time consuming. Such operations also can require the movement of the old roofing materials and objects such as air conditioning units."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of soybean breeding. In particular, the invention relates to the novel soybean variety 01051076.\n2. Description of Related Art\nThere are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include higher seed yield, resistance to diseases and insects, better stems and roots, tolerance to drought and heat, better agronomic quality, resistance to herbicides, and improvements in compositional traits.\nSoybean, Glycine max (L.), is a valuable field crop. Thus, a continuing goal of plant breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties.\nThe oil extracted from soybeans is widely used in food products, such as margarine, cooking oil, and salad dressings. Soybean oil is composed of saturated, monounsaturated, and polyunsaturated fatty acids, with a typical composition of 11% palmitic, 4% stearic, 25% oleic, 50% linoleic, and 9% linolenic fatty acid content (“Economic Implications of Modified Soybean Traits Summary Report,” Iowa Soybean Promotion Board & American Soybean Association Special Report 92S, May 1990)."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a turbine blade tip seal, and more specifically to a turbine blade tip seal with tip cooling.\n2. Description of the Related Art Including Information Disclosed Under 37 CFR 1.97 and 1.98\nIn a gas turbine engine, especially an industrial gas turbine engine, the turbine includes stages of turbine blades that rotate within a shroud that forms a gap between the rotating blade tip and the stationary shroud. Engine performance and blade tip life can be increased by minimizing the gap so that less hot gas flow leakage occurs.\nHigh temperature turbine blade tip section heat load is a function of the blade tip leakage flow. A high leakage flow will induce a high heat load onto the blade tip section. Thus, blade tip section sealing and cooling have to be addressed as a single problem. A prior art turbine blade tip design is shown in FIGS. 1-3 and includes a squealer tip rail that extends around the perimeter of the airfoil flush with the airfoil wall to form an inner squealer pocket. The main purpose of incorporating the squealer tip in a blade design is to reduce the blade tip leakage and also to provide for improved rubbing capability for the blade. The narrow tip rail provides for a small surface area to rub up against the inner surface of the shroud that forms the tip gap. Thus, less friction and less heat are developed when the tip rubs.\nTraditionally, blade tip cooling is accomplished by drilling holes into the upper extremes of the serpentine coolant passages formed within the body of the blade from both the pressure and suction surfaces near the blade tip edge and the top surface of the squealer cavity. In general, film cooling holes are built in along the airfoil pressure side and suction side tip sections and extend from the leading edge to the trailing edge to provide edge cooling for the blade squealer tip. Also, convective cooling holes also built in along the tip rail at the inner portion of the squealer pocket provide additional cooling for the squealer tip rail. Since the blade tip region is subject to severe secondary flow field, this requires a large number of film cooling holes that requires more cooling flow for cooling the blade tip periphery. FIG. 1 shows the prior art squealer tip cooling arrangement and the secondary hot gas flow migration around the blade tip section. FIG. 2 shows a profile view of the pressure side and FIG. 3 shows the suction side each with tip peripheral cooling holes for the prior art turbine blade of FIG. 1.\nThe blade squealer tip rail is subject to heating from three exposed side: 1) heat load from the airfoil hot gas side surface of the tip rail, 2) heat load from the top portion of the tip rail, and 3) heat load from the back side of the tip rail. Cooling of the squealer tip rail by means of discharge row of film cooling holes along the blade pressure side and suction peripheral and conduction through the base region of the squealer pocket becomes insufficient. This is primarily due to the combination of squealer pocket geometry and the interaction of hot gas secondary flow mixing. The effectiveness induced by the pressure film cooling and tip section convective cooling holes become very limited.\nOne prior art reference, U.S. Pat. No. 5,476,364 issued to Kildea on Dec. 19, 1995 and entitled TIP SEAL AND ANTI-CONTAMINATION FOR TURBINE BLADES shows a turbine blade with a tip region having a plurality of spaced holes or slots extending chordwise from the leading edge to the trailing edge of the blade tip, the holes being connected to an internal cooling passage, the holes opening into a cavity formed on the outer surface of the tip. The cavities are separate from each other along the tip edge and all are connected to the internal cooling air passage to discharge cooling air out onto the corner of the tip on the pressure side. The cavities in the Kildea patent are not diffusion openings and do not allow for a formation of a layer of film cooling air on the side of the tip that flows over the tip edge as does the trenches in the present invention.\nThis problem associated with turbine airfoil tip edge cooling can be minimized by incorporation of a new and effective blade tip cooling geometry design of the present invention into the prior art airfoil tip section cooling design."}
{"text": "1. Field\nApparatuses and methods consistent with exemplary embodiments relate to a display apparatus and a method for driving the display apparatus, and more particularly, to a display apparatus of which backlight is operated as being divided into an upper part and a lower part and a method for driving the display apparatus.\n2. Description of the Prior Art\nA liquid crystal display (LCD) apparatus may display a three-dimensional (3D) stereoscopic image as a 120 Hz image. In addition, an LCD apparatus may add a blank image (B) between a left eye image (L) and a right eye image (R) (that is, display an image frame in the order of LBRB) to prevent cross talk effect of the left eye image and the right eye image of the 3D image.\nA backlight of an LCD display apparatus operates as being divided into an upper part and a lower part, and the time when the upper part and the lower part are turned on is different from each other. If the upper part of a backlight is turned on while the lower part of the backlight is turned off or vice versa, brightness in a vertical direction of a screen of the LCD apparatus becomes out of balance.\nA user desires to watch an even and high quality image. Therefore, a method for improving balance in a vertical brightness of an LCD apparatus is required."}
{"text": "1. Field of the Invention\nThe present invention relates to a electrostatically-driven/capacitance-detection type gyro sensor\n2. Description of Related Art\nAs one type of the gyro sensor used for sensing angular velocity, there is known a electrostatically-driven/capacitance-detection type gyro sensor. As shown in FIG. 17, the gyro sensor of this type is constituted by a sensing element 101 including a movable part displaceable in a certain plane (referred to as “motion plane” hereinafter), and a detector unit 102 driving the movable part for detecting the angular velocity of rotative motion imparted to the sensing element around an axis orthogonal to the motion plane (refer to Published Japanese Translation No. 2001-515201 of a PCT Application, for example).\nThe sensing element 101 includes a drive electrode generating an electrostatic force with the movable part to vibrate the movable part along a predetermined direction (refereed to as “driving direction” hereinafter) in the motion plane, a monitor electrode forming a monitoring variable capacitor with the movable part whose capacitance varies depending on the displacement of the movable part along the driving direction, a sense electrode forming a sensing variable capacitor with the movable part whose capacitance varies depending on the displacement of the movable part in the motion plane along a direction orthogonal to the driving direction, and a movable electrode through which the movable part is applied with a bias voltage.\nBy properly controlling the voltage applied between the movable electrode and the drive electrode, it is possible that the movable part continues to vibrate along the driving direction. If the sensing element 101 is imparted with rotative motion around an axis orthogonal to the motion plane when the movable part is vibrating along the driving direction, the movable part develops vibration in the direction orthogonal to the driving direction by the action of the Coriolis force which depends on the angular velocity of the sensing element 101.\nThe detector unit 102 includes a reference voltage generating section 108 generating a reference voltage Vref, a high voltage generating section 106 amplifying the reference voltage Vref to generate the bias voltage VK (=kRK·Vref, kRK being the gain of the high voltage generating section 106) to be applied to the movable electrode, a drive buffer 105 generating a sinusoidal drive signal VD(t) having an offset voltage proportional to the reference voltage Vref, a CV (Capacitance to Voltage) converting section 103 generating a monitor signal VM(t) whose voltage level varies depending on the variation of the capacitance of the monitoring variable capacitor (that is, the vibration state of the movable part vibrating under the action of the drive signal), and generating a sense signal VS(t) whose voltage level varies depending on the variation of the capacitance of the sensing variable capacitor (that is, the vibration state of the movable part vibrating under the action of the Coriolis force), a drive signal controlling section 104 controlling the amplitude and phase of the drive signal VD(t) generated by the drive buffer 105 in accordance with the monitor signal generated by the CV converting section 103 such that the vibration of the movable part along the driving direction continues, and a sensor output signal generating section 107 outputting a sensor output signal VYAW that has an offset voltage proportional to the reference voltage Vref and has a voltage level varying depending on the Coriolis force applied to the movable part (that is, the angular velocity of the movable part).\nIn the gyro sensor having the configuration described above, the drive force driving the movable part depends on the voltage difference between the movable electrode and the drive electrode, and the conversion gain of the CV converting section 103 generating the monitor signal VM(t) and sense signal VS(t) depends on the voltage differences between the movable electrode and the monitor electrode or sense electrode. On the other hand, the gain for converting the rotative motion imparted to the sensing element 101 into the Coriolis force (referred to as “element sensitivity” hereinafter) depends on the vibrating state of the movable part vibrating under the action of the drive force.\nAccordingly, if the reference voltage Vref varies depending on the variation of the power supply voltage VCC, the gain for converting the angular velocity of the sensing element 101 into the sensor output signal VYAW (referred to as “sensor sensitivity” hereinafter) changes, since the element sensitivity and the circuit characteristics of the detector unit 102 including the CV converting section 103 change depending on the variation of the power supply voltage VCC.\nAccordingly, in order that the sensor output signal VYAW is not affected by the variation of the power supply voltage VCC, the reference voltage generating section 108 is configured to generate the reference voltage Vref by use of a constant voltage not affected by the variation of the power supply voltage VCC, such as a band-gap voltage.\nIn a microcomputer-based system using such a gyro sensor, it is common that the sensor output signal VYAW outputted from the gyro sensor is converted to a digital signal by use of an A/D converter, and then supplied to a microcomputer.\nGenerally, in such a system, a reference voltage which the A/D converter uses for performing the A/D conversion is generated by dividing down the power supply voltage VCC. Since the sensor output signal VYAW varies between both polarities with respect to a predetermined offset voltage, the microcomputer processes the digital signal outputted from the A/D converter by using the center value of the output range of the A/D converter as a zero point. Generally, in order to maximize the dynamic range of the sensor output signal VYAW, the offset voltage corresponding to the zero point is set at half the power supply voltage VCC.\nAccordingly, when the power supply voltage VCC is 5V, the offset voltage thereof is set at 2.5V. As explained above, although the zero point of the sensor output signal VYAW outputted from the gyro sensor is set in accordance with the reference voltage Vref, and is therefore maintained unchanged even if the power supply voltage VCC varies, the zero point recognized by the microcomputer varies in proportion to the power supply voltage VCC. Accordingly, if the power supply voltage VCC deviates from its rated value for some reason, there arises a difference between the zero point recognized by the microcomputer and the zero point of the sensor output signal VYAW outputted from the gyro sensor. This lowers the accuracy of angular velocity detection.\nFurthermore, such a system has another problem in that, if the power supply voltage VCC deviates from its rated value, although there occurs no change in the sensor output signal VYAW, there occurs a change in the dynamic range of the A/D converter. This causes a change in the incremental value of the sensor output signal VYAW per bit of the digital signal which the A/D converter generates. This changes the conversion ratio of the A/D converter, which also lowers the accuracy of angular velocity detection.\nIt may occur that the accuracy of the angular velocity detection can be avoided from being lowered if the microcomputer operates to detect the deviation of the power supply voltage VCC or the deviation of the reference voltage Vref due to the deviation of the power supply voltage VCC, and to compensate for the deviation of the zero point or the sensitivity in accordance with the result of the detection. However, this considerably increases the burden of the microcomputer."}
{"text": "1. Field of the Invention\nThe invention relates to a transparent composite system comprising polymethyl methacrylate (PMMA) plates and an intermediate layer of an elastomeric resin.\n2. Discussion of the Background\nPolymethyl methacrylate, and copolymers of methyl methacrylate, exhibit characteristics which are desirable for glazing materials: particularly high transparency and good weather resistance, and--in comparison with silicate glass--a lower weight. But these advantageous, optical characteristics of PMMA and its copolymers are diminished by the brittleness of these materials. This brittleness exhibits itself in the splintering off of shards or parts of the PMMA when it undergoes mechanical stress.\nThis intrinsic brittleness of PMMA (in addition to its poor scratch resistance and its flammability) naturally reduces its application, especially as a construction material or in areas of high traffic or public use. Several attempts to eliminate these disadvantages have been reported. (see H. Rauch-Puntigam, Th. Volker, Acryl- und Methacrylverbindungen, Springer-Verlag 1967, p. 291).\nFor example, German Patent Application 29 14 223 recommends, e.g. a clear, transparent synthetic resin composite plate consisting of a clear, transparent plastic base plate with a thickness of at least 1 mm and with parallel surfaces, at least one clear, transparent plastic foil with a scratch resistance not less than that of the base plate, and at least one clear, transparent adhesive resin layer positioned between the base plate and plastic foil, which adheres to the foil in fixed manner and which can be pulled off the base plate.\nJapanese Application 61 132 343 (Chem. Abstr. 105, 154387e) discloses an acrylic glass laminate consisting of two exterior hard plastic plates and an intermediate plate of soft plastic with an electrically conductive layer on or near the surface. The intermediate layer is, e.g., an MMA-methacrylate-elastomer equipped on both sides with an antistatic agent and an acrylate adhesive.\nGerman Application 31 28 985 discloses laminated, thermoplastic objects. The laminates are obtained by joining at least one PMMA film with at least one impact-resistant methacrylate polymer by calendaring.\nEuropean Patent 419 166 describes co-extruded multi-layer webs consisting of a complex configuration of individual layers with a thickness of 0.3 to 6 mm each. The first exterior layer consists of a blend of polyvinylidene fluoride with M.sub.n &lt;250,000, 30-60% by weight C1-C8 alkyl methacrylate with M.sub.n .ltoreq.150,000, and possibly 10-30% by weight of a graft elastomer of C1-C8 alkyl acrylate, grafted and cross-linked with allylmethacrylate, as well as at least one second layer of an alkyl methacrylate homopolymer or copolymer. The co-extrusion is performed through a multiple extrusion nozzle.\nThere are also papers on laminates from different plastics. Illinger et al. in Polym. Sci. Technol. 1975, 9A (Adhes. Sci. Technol) 217-232, e.g., proposes laminates of acrylate-polyurethane-polycarbonate for the protection of windshields, whereby it is emphasized that the impact resistance and the optical clarity of the laminates depend on the chemical composition of the intermediate adhesion layer. European Patent Application 138 083 proposes polymer mixtures as an intermediate adhesion layer for laminates, said intermediate adhesion layers consisting of copolymers of firstly 81-99 weight parts olefin and unsaturated carboxylic acids and secondly of 1-19 weight parts thermoplastic elastomers, whereby the former are present dispersed as particles. The use of fiber inserts, particularly glass and carbon fiber inserts, has also been considered in attempting to improve the splintering safety of thermoplasts, such as acrylic glasses (cf. European Patent 0 510 927)."}
{"text": "Phosphatidylinositol (hereinafter abbreviated as “PI”) is one of the phospholipids in cell membranes. In recent years it has become clear that PI plays an important role also in intracellular signal transduction. It is well recognized in the art that PI (4,5) bisphosphate (PI(4,5)P2 or PIP2) is degraded into diacylglycerol and inositol (1,4,5) triphosphate by phospholipase C to induce activation of protein kinase C and intracellular calcium mobilization, respectively [M. J. Berridge et al., Nature, 312, 315 (1984); Y. Nishizuka, Science, 225, 1365 (1984)].\nPhosphatidylinositol-3 kinase (“PI3K”) is an enzyme that phosphorylates the 3-position of the inositol ring of phosphatidylinositol [D. Whitman et al., Nature, 332, 664 (1988)]. Pluralities of PI3K subtypes exist. Three major subtypes of PI3Ks have now been identified on the basis of their in vitro substrate specificity, and these three are designated class 1 (a & b), class II, and class III [B. Vanhaesebroeck, Trend in Biol. Sci., 22, 267 (1997)].\nThe class Ia PI3K subtype has been most extensively investigated to date. Within the class Ia subtype there are three isoforms (α, β, & δ) that exist as hetero dimers of a catalytic 110-kDa subunit and regulatory subunits of 50-85 kDa. The regulatory subunits contain SH2 domains that bind to phosphorylated tyrosine residues within growth factor receptors or adaptor molecules and thereby localize PI3K to the inner cell membrane. At the inner cell membrane PI3K converts PIP2 to PIPS (phosphatidylinositol-3,4,5-trisphosphate) that serves to localize the downstream effectors PDK1 and Akt to the inner cell membrane where Akt activation occurs. Activated Akt mediates a diverse array of effects including inhibition of apoptosis, cell cycle progression, response to insulin signaling, and cell proliferation. Class Ia PI3K subtypes also contain Ras binding domains (RBD) that allow association with activated Ras providing another mechanism for PI3K membrane localization. Activated, oncogenic forms of growth factor receptors, Ras, and even PI3K kinase have been shown to aberrantly elevate signaling in the PI3K/Akt/mTOR pathway resulting in cell transformation. As a central component of the PI3K/Akt/mTOR signaling pathway PI3K (particularly the class Ia a isoform) has become a major therapeutic target in cancer drug discovery.\nSubstrates for class 1 PI3Ks are PI, PI(4)P and PI(4,5)P2, with PI(4,5)P2 being the most favored. Class 1 PI3Ks are further divided into two groups, class Ia and class Ib, because of their activation mechanism and associated regulatory subunits. The class Ib PI3K is p110γ that is activated by interaction with G protein-coupled receptors. Interaction between p110γ and G protein-coupled receptors is mediated by regulatory subunits of 110, 87, and 84 kDa.\nPI and PI(4)P are the known substrates for class II PI3Ks; PI(4,5)P2 is not a substrate for the enzymes of this class. Class II PI3Ks include PI3K C2α, C2β and C2γ isoforms, which contain C2 domains at the C terminus, implying that their activity is regulated by calcium ions.\nThe substrate for class III PI3Ks is PI only. A mechanism for activation of the class III PI3Ks has not been clarified. Because each subtype has its own mechanism for regulating activity, it is likely that activation mechanism(s) depend on stimuli specific to each respective class of PI3K.\nThe compound PI103 (3-(4-(4-morpholinyl)pyrido[3′,2′:4,5]furo[3,2-d]pyrimidin-2-yl)phenol) inhibits PI3Ka and P13 Kg as well as the mTOR enzymes with IC50 values of 2, 3, and 50-80 nM respectively. I.P. dosing in mice of this compound in human tumor xenograft models of cancer demonstrated activity against a number of human tumor models, including the glioblastoma (PTEN null U87MG), prostate (PC3), breast (MDA-MB-468 and MDA-MB-435) colon carcinoma (HCT 116); and ovarian carcinoma (SKOV3 and IGROV-1); (Raynaud et al, Pharmacologic Characterization of a Potent Inhibitor of Class I Phosphatidylinositide 3-Kinases, Cancer Res. 2007 67: 5840-5850).\nThe compound ZSTK474 (2-(2-difluoromethylbenzoimidazol-1-yl)-4,6-dimorpholino-1,3,5-triazine) inhibits PI3Ka and P13 Kg but not the mTOR enzymes with an IC50 values of 16, 4.6 and >10,000 nM respectively (Dexin Kong and Takao Yamori, ZSTK474 is an ATP-competitive inhibitor of class I phosphatidylinositol 3 kinase isoforms, Cancer Science, 2007, 98:10 1638-1642). Chronic oral administration of ZSTK474 in mouse human xenograft cancer models, completely inhibited growth which originated from a non-small-cell lung cancer (A549), a prostate cancer (PC-3), and a colon cancer (WiDr) at a dose of 400 mg/kg. (Yaguchi et al, Antitumor Activity of ZSTK474, a New Phosphatidylinositol 3-Kinase Inhibitor, J. Natl. Cancer Inst. 98: 545-556).\nThe compound NVP-BEZ-235 (2-methyl-2-(4-(3-methyl-2-oxo-8-(quinolin-3-yl)-2,3-dihydro-1H-imidazo[4,5-c]quinolin-1-yl)phenyl)propanenitrile) inhibits both PI3Ka and PI3Kg as well as the mTOR enzymes with IC50 values 4, 5, and “nanomolar”. Testing in human tumor xenograft models of cancer demonstrated activity against human tumor models of prostrate (PC-3) and glioblastoma (U-87) cancer. It entered clinical trials in December of 2006 (Verheijen, J. C. and Zask, A., Phosphatidylinositol 3-kinase (PI3K) inhibitors as anticancer drugs, Drugs Fut. 2007, 32(6): 537-547).\nThe compound SF-1126 (a prodrug form of LY-294002, which is 2-(4-morpholinyl)-8-phenyl-4H-1-benzopyran-4-one) is “a pan-PI3K inhibitor”. It is active in preclinical mouse cancer models of prostrate, breast, ovarian, lung, multiple myeloma, and brain cancers. (Verheijen, J. C. and Zask, A., Phosphatidylinositol 3-kinase (PI3K) inhibitors as anticancer drugs, Drugs Fut. 2007, 32(6): 537-547).\nAlthough it seems clear that inhibition of the α isoform is essential for the antitumor activity of PI3K inhibitors, it is not clear whether a more selective inhibitor of a particular PI3K isoform may lead to fewer unwanted biological effects. It has recently been reported that non-PI3Kα class I isoforms (PI3Kβ, δ and γ) have the ability to induce oncogenic transformation of cells, suggesting that nonisoform-specific inhibitors may offer enhanced therapeutic potential over specific inhibitors.\nSelectivity versus other related kinases is also an important consideration for the development of PI3K inhibitors. While selective inhibitors may be preferred in order to avoid unwanted side effects, there have been reports that inhibition of multiple targets in the PI3K/Akt pathway (e.g., PI3Kα and mTOR [mammalian target of rapamycin]) may lead to greater efficacy. It is possible that lipid kinase inhibitors may parallel protein kinase inhibitors in that nonselective inhibitors may also be brought forward to the clinic.\nMammalian Target of Rapamycin, mTOR, is a cell-signaling protein that regulates the response of tumor cells to nutrients and growth factors, as well as controlling tumor blood supply through effects on Vascular Endothelial Growth Factor, VEGF. Inhibitors of mTOR starve cancer cells and shrink tumors by inhibiting the effect of mTOR. All mTOR inhibitors bind to the mTOR kinase. This has at least two important effects. First, mTOR is a downstream mediator of the PI3K/Akt pathway. The PI3K/Akt pathway is thought to be over activated in numerous cancers and may account for the widespread response from various cancers to mTOR inhibitors. The over-activation of the upstream pathway would normally cause mTOR kinase to be over activated as well. However, in the presence of mTOR inhibitors, this process is blocked. The blocking effect prevents mTOR from signaling to downstream pathways that control cell growth. Over-activation of the PI3K/Akt kinase pathway is frequently associated with mutations in the PTEN gene, which is common in many cancers and may help predict what tumors will respond to mTOR inhibitors. The second major effect of mTOR inhibition is anti-angiogenesis, via the lowering of VEGF levels.\nIn lab tests, certain chemotherapy agents were found to be more effective in the presence of mTOR inhibitors. George, J. N., et al., Cancer Research, 61, 1527-1532, 2001. Additional lab results have shown that some rhabdomyosarcoma cells die in the presence of mTOR inhibitors.\nThere are three mTOR inhibitors, which have progressed into clinical trials. These compounds are Wyeth's Torisel, also known as 42-(3-hydroxy-2-(hydroxymethyl)-rapamycin 2-methylpropanoate, CCI-779 or Temsirolimus; Novartis' Everolimus, also known as 42-O-(2-hydroxyethyl)-rapamycin, or RAD 001; and Ariad's AP23573 also known as 42-(dimethylphopsinoyl)-rapamycin. The FDA has approved Torisel for the treatment of advanced renal cell carcinoma. In addition, Torisel is active in a NOS/SCID xenograft mouse model of acute lymphoblastic leukemia [Teachey et al, Blood, 107(3), 1149-1155, 2006]. On Mar. 30, 2009, the U.S. Food and Drug Administration (FDA) approved Everolimus (AFINITOR™) for the treatment of patients with advanced renal cell carcinoma. AP23573 has been given orphan drug and fast-track status by the FDA for treatment of soft-tissue and bone sarcomas.\nThe three mTOR inhibitors have non-linear, although reproducible pharmacokinetic profiles. Mean area under the curve (AUC) values for these drugs increase at a less than dose related way. The three compounds are all semi-synthetic derivatives of the natural macrolide antibiotic rapamycin. It would be desirable to find fully synthetic compounds, which inhibit mTOR that are more potent and exhibit improved pharmacokinetic behaviors."}
{"text": "Various kinds of nitride semiconductor devices using a nitride semiconductor, such as gallium nitride, have been developed. Examples thereof include semiconductor light emitting devices, such as an ultraviolet, blue, or green light emitting diode (LED), and a purple-blue or blue laser diode (LD).\nA GaN layer of the nitride semiconductor device is usually formed on a sapphire substrate or the like. However, due to a reduction in cost and an increase in diameter of a wafer, it is desirable to establish a technique for growing the GaN layer on a substrate other than the sapphire substrate."}
{"text": "The present invent relates to a camera, and more particularly to a camera which can be mounted in box and dome camera configurations.\nCameras, also referred as camera lenses, are widely used in the area of field monitoring, video recording, conducting surveillance and video recording of traffic conditions, production status, and status of a particular area.\nCurrently, cameras are generally utilized in two configurations, namely, the box camera configuration and the dome camera configuration. In the box camera configuration, a camera is substantially housed in a columnar housing with a lens extending out from the front end of the housing, and the housing of the camera is mounted on, for example, a wall or a traffic signal mount by using a mounting bracket, the lens being positioned at a particular viewing angle. In addition, under desired circumstances, said bracket may have a function of panning or tilting, or both, to adjust the viewing angle of the camera. Such a configuration is commonly seen in the outdoor mounting of cameras. In the dome camera configuration, a camera is generally suspended and mounted onto the ceiling, the entirety of which is covered by a hemispherical transparent dome, and outer images can be captured by a lens of the camera through the hemispherical transparent dome. This lens is generally supported by a gimbal member, and the viewing angle is controlled by an appropriate control mechanism. The camera in dome camera configuration has relatively nice looking and occupies a relatively small volume; therefore, dome camera configurations are commonly used indoors or in an elevator."}
{"text": "The invention relates to an electrical terminal connection for the electrodes of a gas discharge over or surge voltage arrester with the following features:\n(a) the electrodes to be connected consist of a material with good current conductivity;\n(b) the connection should respectively occur at a wire consisting of a material with good current conductivity.\nEfficient over-voltage arresters must be able to arrest surge currents&gt;10 kA. In dual-path over-voltage arresters, this meas that a sum current&gt;20 kA flows across a common center electrode. Given electrical terminal connections via resiliently contacting contacts, such high surge currents cause sputtering of the contact material and thereby damage the sockets. Terminal wires which are soldered in and which generally consist of copper and then exhibit a diameter of approximately 1 mm offer a better electrical terminal connection for such currents. Copper is likewise preferably selected as the electrode material for electrical and economic reasons. The rigid connection between the electrode and the terminal wire can then be produced by means of soldering or welding.\nBoth are proposed, for example, in German OS No. 28 28 650, incorporated herein by reference. On the one hand, the possibility is there described of soldering a tubular rivet to the floor of the cup-shaped electrode and then squeezing (FIG. 4) the terminal wire in the tubular rivet. On the other hand, a welded connection is provided to \"press\" a cylinder to the floor of the cup-shaped electrode, the diameter of said cylinder being approximately 1.5 times as great as the diameter of the terminal wire. The terminal wire is welded blunt to said cylinder. Said cylinder has the object of reducing the thermal dissipation during the welding operation. Welding to the smooth bottom of the electrode or, respectively, to the jacket surface of a center electrode would not be possible given material which has high current conductivity and, thus also has high thermal conductivity.\nEven given a cylinder dimensioned in such manner, however, a welded connection is only possible as a short-time welding, for instance as pulsed resistance welding. Particularly given a connection of copper to copper, this, however produces unreliable connections."}
{"text": "A pervasive trend in modern integrated circuit manufacture is to produce semiconductor devices, such as memory cells, that are as small as possible. A typical memory cell, which generally is formed from a field effect transistor (FET), includes a source and a drain formed in an active region of a semiconductor substrate by implanting N-type or P-type impurities in the semiconductor substrate. Disposed between the source and the drain is a channel (or body) region. Disposed above the body region is a gate electrode. The gate electrode and the body are spaced apart by a gate dielectric layer. It is noted that memory cells can be formed in bulk format (for example, the active region being formed in a silicon substrate) or in a semiconductor-on-insulator (SOI) format (for example, in a silicon film that is disposed on an insulating layer that is, in turn, disposed on a silicon substrate).\nAlthough the fabrication of smaller transistors allows more transistors to be placed on a single monolithic substrate for the formation of relatively large circuit systems in a relatively small die area, this downscaling can result in a number of performance degrading effects. In FET devices with a channel having a relatively short length, the FET can experience a number of undesirable electrical characteristics referred to as short channel effect (SCE). SCE generally occur when the gate does not have adequate control over the channel region, and can include threshold voltage (Vt) roll-off, off current (Ioff) roll-up and drain induced barrier lowering (DIBL). As the physical dimensions decrease, SCE can become more severe. SCE is the result of intrinsic properties of the crystalline materials used in the FET devices. Namely, the band gap and built-in potential at the source/body and drain/body junctions are non-scalable with the reduction of physical device dimensions, such as a reduction in channel length.\nAnother problem due to downscaling of transistors is charge loss in the oxide-nitride-oxide (ONO) layer. Charge loss is due in part to ONO damage during an erase cycle. During the erase cycle, the relatively high voltage (Vpp) applied across the oxide causes tunneling of electrons from the charge trapping nitride layer to the drain region. At the same time, the high voltage could cause holes from the drain to be injected into the oxide. These holes can degrade the performance of the oxide by creating a leakage path in the oxide between the drain and the nitride layer.\nSince the oxide is the barrier for electrons traveling to and from the charge trapping layer, the charging and discharging current of a memory cell depends on the voltage applied across the oxide layer, I=C(dv/dt). Therefore, the voltage applied across the oxide has a direct effect on electron tunneling and is the main cause of undesirable hole injection into the oxide during an erase operation.\nTherefore, there exists a need in the art for semiconductor devices, such as memory cells, that have reduced SCE and for fabrication techniques to make those semiconductor devices. Furthermore, there is a need in the art for semiconductor devices that have reduced ONO damage during erase operations of the device."}
{"text": "1. Field of the Invention\nThis invention relates to a method of coalescing laser beams, and more particularly to a method of coalescing two or more laser beams on an object to get a multiplied power of the laser beams on the object. The \"coalescing\" means mingling or doubling of two or more laser beams on an object onto which the laser beams are to impinge in the doubled or multiplied power.\nIn the present invention, the two or more laser beams are first merged into a single laser beam in which the laser beams are substantially aligned with each other or one another to make a doubled or multiplied laser beam, in which however the laser beams are usually not perfectly merged into a single laser beam. In this specification, however, the laser beams thus merged or roughly mixed together are called \"merged laser beams\" or \"a single merged laser beam\" because they are substantially \"merged\". In a more precise sense, on the other hand, these laser beams first merged into a substantial single laser beam are not perfectly aligned with each other or one another but are out of perfect parallelism and out of perfect coincidence or alignment. In this sense, the laser beams thus merged are then \"collimated\" in a precise sense so that they are directed to a common focal point through a converging lens. Of course, the collimated laser beams after \"merged\" are so close to each other or one another that they advance through the substantially the same part of the same optical elements such as lenses and mirrors.\n2. Description of the Prior Art\nLaser beams are widely used for optical communication, measurement, processing such as laser welding, laser drilling, and the like. When two or more laser beams are used, it is often required to make them coalesced to get a multiplied power or a laser beam of wider wavelength coverage.\nFor example, in printing scanners, when a color image original is read out by use of a He-Ne laser beam, an Ar.sup.+ laser beam and a He-Cd laser beam and a monochromatic image is reproduced on a photosensitive material by use of the electric signal obtained by the read-out, it is necessary to converge three laser beams to a single point on the color image original for conducting the read-out. Also, in a radiation image recording and reproducing system as disclosed, for example, in U.S. Pat. Nos. 4,258,264 and 4,346,295, and Japanese Unexamined Patent Publication No. 56(1981)-11395, a stimulable phosphor sheet carrying a radiation image stored therein is scanned with a laser beam which causes the stimulable phosphor sheet to emit light in proportion to the radiation energy stored. The emitted light is detected and converted into an electric signal, and a photosensitive material is scanned with a laser beam modulated based on the electric signal by a light modulator, thereby reproducing an image in the photosensitive material. In this case, in order to increase the read-out speed and the reproducing speed, it is necessary to increase the scanning speed of the laser beam. However, if only the scanning speed of the laser beam is increased with the power thereof maintained on the same level, the intensity of the laser beam per unit area of the surface of the stimulable phosphor sheet or a film scanned thereby inevitably becomes low. Therefore, it is necessary to increase the scanning speed and, at the same time, to increase the power of the laser beam. In order to increase the power of the laser beam, the output power of the laser beam source should be increased. However, when the output power of the laser beam source is increased, the size of the laser beam source inevitably becomes large, and the problem that the laser beam source cannot be incorporated in the radiation image recording and reproducing system or the problem that the radiation image recording and reproducing system must be made large arises. One approach to elimination of the problem is to use two laser beam sources and to coalesce the laser beams emitted thereby by use of an optical merging means such as a polarization beam splitter or a diffraction grating, thereby multiplying the laser beam power. This method is based on the finding that, in the case where the same level of laser beam power should be obtained, it is possible to make the size of the laser beam emitting system smaller when two laser beam sources are used and the laser beams emitted thereby are merged than when a single laser beam source is used. However, also in this method, it is required to make the two laser beams align with each other. The same thing is required also in laser processing apparatuses such as laser welding apparatuses.\nHowever, the step of making two or more very thin laser beams align with each other is not always easy to conduct, or requires a complicated mechanism, or requires much time to adjust the mechanism."}
{"text": "We start by presenting some basic concepts to facilitate the understanding of the numerous policing techniques that are presented afterwards.\nPackets\nA data sender usually splits data to be sent into small units known as packets. Each packet consists of a header and a payload carrying the data to be delivered. The header contains fields defined by the relevant communication protocol. The great majority of packets carried by commercial networks nowadays are so-called IP packets. IP is the Internet Protocol. This ensures that a network of routers can forward any packet from the source to its destination. IP is a connectionless protocol—that means that the header information in each data packet is sufficiently self-contained for routers to deliver it independently of other packets; each packet could even take a different route to reach the destination.\nDistributed Bandwidth Sharing and Congestion\nData traversing the Internet follows a path between a series of routers, controlled by various routing protocols. Each router seeks to move the packet closer to its final destination. If too much traffic traverses the same router in the network, the router can become congested and packets start to experience excessive delays whilst using that network path. If sources persist in sending traffic through that router it could become seriously overloaded (congested) and even drop traffic (when its buffers overflow). If sources still persist in sending traffic through this bottleneck it could force more routers to become congested, and if the phenomenon keeps spreading, that can lead to a congestion collapse for the whole Internet—which occurred regularly in the mid-1980s.\nThe solution to that problem has been to ensure that sources take responsibility for the rate at which they send data over the Internet by implementing congestion control mechanisms. Sources monitor feedback from the receiver of the metric that characterises path congestion in order to detect when the path their data is following is getting congested, in which case they react by reducing their throughput—while they may slowly increase their rate when there is no sign of the path becoming congested.\nTypical path characterisation metrics that sources monitor are average roundtrip time (RTT) for the data path, variance of the roundtrip time (jitter) and level of congestion on the path.\nThe congestion level can be signalled either implicitly (through congested routers dropping packets when their buffers overflow or to protect themselves) or explicitly (through mechanisms such as explicit congestion notification—see next subsection). Currently the most common option is implicit signalling.\nSources using TCP are able to detect losses, because a packet loss causes a gap in the sequence; whenever a TCP source detects a loss, it is meant to halve its data transmission rate, but no more than once per round trip time, which alleviates the congestion on the router at the bottleneck.\nDEC-Bit Scheme\nThe DEC-bit scheme was the ancestor of modern mechanisms to convey congestion notification in packet networks by providing explicit congestion notification to the sources. This is discussed in a paper entitled “A Binary Feedback Scheme for Congestion Avoidance in Computer Networks” by Ramakrishnan and Jain, which will be referred to for convenience as [RAN90a]. Bibliographic details of this and other prior art documents are provided in the “References” section below. With this scheme when a router detects congestion it sets a bit (CI) in each packet, the receiver then relays this information back to the sender in its acknowledgments, and the sender in turn adjusts its transmission rate. In the router the congestion detection algorithm is based on average queue length. The source reaction to DEC-bit in acknowledgments is an Additive Increase Multiplicative Decrease (AIMD) response which means that congestion window grows at linear rate in the absence of congestion feedback and it declines multiplicatively (i.e. exponentially) for every congestion feedback event.\nExplicit Congestion Notification\nExplicit Congestion Notification (ECN) [RFC3168] conveys congestion in TCP/IP networks by means of two-bit ECN field in the IP header, whether in IPv4 (see FIG. 1) and IPv6 (see FIG. 2). Prior to the introduction of ECN, these two bits were present in both types of IP header, but always set to zero. Therefore, if these bits are both zero, a queue management process assumes that the packet comes from a transport protocol on the end-systems that will not understand the ECN protocol so it only uses drop, not ECN, to signal congestion.\nThe meaning of all four combinations of the two ECN bits is shown in FIG. 3. If either bit is one, it tells a queue management process that the packet has come from an ECN-capable transport (ECT) that will understand ECN marking, as well as drop, as a signal of congestion. When a queue management process detects congestion, for packets with a non-zero ECN field, it sets the ECN field to the Congestion Experienced (CE) codepoint. On receipt of such a marked packet, a TCP receiver sets the Echo Congestion Experienced (ECE) flag in the TCP header, which the TCP source interprets as if the packet has been dropped for the purpose of its rate control.\nDrop and congestion signals are not mutually exclusive signals, and flows that enable ECN have the potential to detect and respond to both signals.\nAs well as the idea of ECN being adopted in IP, it was previously adopted in Frame Relay and ATM, but in these latter two protocols the network arranges feedback of the congestion signals internally, and the network enforces traffic limits to prevent congestion build-up [ITU-T Rec.I.371]. When a cell arriving at a switch causes a congestion threshold to be exceeded, then its EFCI bit is set.\nThe IEEE has standardised an explicit congestion approach where Ethernet switches not the end systems arrange to feedback the congestion signals, although the Ethernet device on the sending system is expected to co-operate by reducing its rate in response to the signals. The approach is tailored exclusively for homogeneous data centre environments.\nIn previous approaches, each frame (or packet) carried just a binary flag and the strength of the congestion signal depended on the proportion of marked frames—effectively a unary encoding of the congestion signal in a stream of zeroes and ones. However, the IEEE scheme signals a multi-bit level of congestion in each feedback frame, hence its name: Quantised Congestion Notification or QCN [IEEE802.1Qau].\nRandom Early Detection (RED)\nHistorically, routers would drop packets when they got completely saturated (i.e. when a traffic burst cannot be accommodated in the buffer of the router)—this policy is called drop-tail. Random early detection (RED) [RED] is an improvement intended to desynchronise TCP flows (synchronisation occurs when multiple TCP flows increase and decrease their transmission window at the same time). RED is an active queue management (AQM) process that monitors the average queue length in a buffer and when this is higher than a given threshold, the router starts to drop/mark packets with a probability which increases with the excess length of the averaged queue over the threshold. RED is widely used in today's Internet because it allows sources to react more promptly to incipient congestion and it keeps queues from growing unnecessarily long. Equipment vendors have implemented variants of RED, e.g. Cisco' proprietary implementation is Weighted Random Early Detection (WRED). However, it is well known that RED is very sensitive to parameter setting.\nIn the course of experimental work to investigate the sensitivity of RED to its parameter settings, the timescale over which RED averaged the queue was investigated, and in one case at least averaging was turned off completely, instead using the instantaneous queue [RED-params].\nThe use of an Active Queue Management (AQM) technique using Random Early Detection (RED) will be described with reference to FIG. 4.\nRED randomly drops/marks packets with a probability p that depends upon the smoothed queue qave. In a RED-based AQM, a smoothed queue qave is continuously estimated by means of an exponential weighted moving average (EWMA) of the real queue q:qave←(1−wq)qave+wqq where wq is the weight given to the real queue's length; see FIG. 4a) for an example of real queue versus smoothed queue evolution over time.\nWhen the smoothed queue size qave is below a minimum threshold q0 then no packets are dropped/marked. When qave is between q0 and q1 then packets are discarded with a probability p, between 0 and p1, that is linearly proportional to qave. When qave is greater than threshold q1 then probabilistic drop/mark continues with an increased probability ranging between p1 and pmax, which still depends linearly on qave. This is also known as the Gentle variant of RED algorithm [GRED] (see FIG. 4b)). Its behaviour differs from the original RED between q1 and qmax: while GRED drops/marks packets with a probability which depends linearly upon smoothed queue length with a maximum probability pmax, in such interval RED marks with a maximum probability pmax. See FIG. 4c) for an example of time evolution of packet drop/mark probability. Assuming the Sources of packet load are responsive to congestion, it can be seen that a period of drops/marks p>0 results in senders slowing down, which results in reduced queue length and hence no drop/mark p=0; subsequently senders increase their transmission rate which results in an increased queue length, note that it takes some time before p>0 again because of the queue length smoothing.\nFairness Improvement for RED (FI-RED)\nA paper entitled “FI-RED: AQM Mechanism for Improving Fairness among TCP Connections in Tandem Networks” by H. Ohsaki et al [FI-RED] discusses an active queue management mechanism for improving fairness among heterogeneous TCP connections. In FI-RED, an ECN-based technique is used for differentiating the packet marking probability of RED in dependence on whether the CE (Congestion Experienced) bit is set in arriving packets. With FI-RED, congestion indications to TCP connections with a large number of hops (i.e. those connections with a high probability that the CE bit is set) are suppressed. By virtue of such a modification to RED, a TCP connection with a large number of hops will suffer almost the same packet marking probability as one with a small number of hops. This is intended to improve fairness among TCP connections with different numbers of hops. It will be noted that with FI-RED, ECN marks are assigned to packets with a marking probability based on pre-existing congestion markings on the arriving packets (not the ECN-capability itself), which is then compared with the marking probability determined from the local queue length by the traditional RED algorithm.\nRandom Exponential Marking (REM) [REM]\nWhereas with RED, congestion notification depends upon the average queue length, with the REM AQM, congestion signals depend on rate deviation (i.e. queue growth over a period) and queue deviation (i.e. difference between queue length and a target queue length). The queue is sampled periodically. Between samples, the congestion signal thus depends on a delayed rather than an instantaneous measure of the queue. The probability of marking or dropping in the REM AQM is characterised by an exponential as opposed to a piece-wise linear function used by RED, and similarly to RED congestion feedback can be by dropping or by ECN marking. The exponential function in REM ensures that signals accumulated along a path are precisely summed to signal end-to-end congestion to the source.\nBLUE\nBLUE [BLUE] is an AQM designed to overcome the inefficiencies that arise in the presence of bursty traffic and AQMs based on a smoothed queue such as RED. With bursty traffic, queue length often oscillates wildly and RED-like AQMs that drop based on smoothed queue length are unable to react sufficiently quickly. BLUE maintains a probability of packet marking or dropping which is based on actual queue occupancy hence, congestion notification (drop or ECN-mark) depends upon actual queue length and not smoothed queue length. This marking probability is increased when packets are being dropped and is decreased when there is less congestion.\nProportional-Integral (PI)\nProportional-integral [PI] controllers are controllers based on feedback control theory which perform AQM functions. Like BLUE, a key feature of the PI controller is to use the instantaneous queue size not an averaged queue length to determine marking and/or loss probability, in order to minimise delay of the feedback signal.\nRED in a Different Light (nRED)\nnRED is a self-configuring variant of RED which is available online [nRED] but has not been published in peer reviewed journals or conferences. This AQM estimates the persistent queue, which differs from the smoothed-RED queue as it only seeks to track persistent buffer occupancy that has not cleared within a given time interval. This is achieved by means of a filtered function of the instantaneous queue length. It compares the instantaneous queue to the filtered value and, if it is greater, the filtered function rises using the same exponentially weighted moving average as RED, but if it is less, the filtered value tracks the instantaneous queue downwards.\nControlled Delay (CoDel)\nCoDel [CoDel] superseded nRED. It uses the service time (or sojourn time) that a packet spends in a queue as the characterisation of queue length and produces a smoothed value of recent sojourn times by taking the minimum sojourn time seen over a recent time interval.\nSample Path Shadow Pricing (SPSP)\nSample Path Shadow Pricing (SPSP) [GI99] is a scheme which prescribes to mark packets that cause the unsmoothed queue to exceed a threshold and then lead to buffer overflow. This is unattainable in practice, as a packet may have left the queue by the time a decision is made that it should have been marked. However, it is used as a theoretical ideal e.g. to define the SPSP marking fairness principle, that is, a marking algorithm satisfies SPSP fairness if it marks the same number of packets that SPSP would have marked. SPSP purely concerns ECN marking and does not consider loss.\nPre-Congestion Notification (PCN)\nThe IETF pre-congestion notification (PCN) architecture is a framework that supports Quality of Service (QoS) of inelastic traffic within a DiffServ domain. Admission control (or flow termination) of inelastic traffic in a PCN-domain depends upon the marked PCN-traffic observed within the priority class reserved for PCN-traffic. PCN traffic can coexist with other traffic and receive different treatment within the queue. Two PCN marking algorithms have been standardised[PCN] based on the idea of virtual queue [GI99] which is a fictitious queue that emulates the behaviour of a real queue if it were draining at a slower rate that the actual link capacity.\nThe standardised PCN marking algorithms both use the instantaneous length of the virtual queue, not a smoothed queue length as with DEC-bit or RED. PCN uses the Not-ECT codepoint in the ECN field (FIG. 3) in a similar way to ECN, in order to support a mixture of traffic sharing the same queue, some of which is capable of being PCN-marked and some hot. Non-PCN traffic sharing a priority queue would be dropped rather than marked if it exceeded the policed rate of the priority queue.\nCongestion Control for High Bandwidth-Delay-Product (BDP) Networks\nTraditional TCP/IP congestion control and recovery mechanics can be detrimental to the performance of delay sensitive applications. This is particularly problematic for high Bandwidth-Delay-Product (BDP) networks (also known as Long Fat Networks (LFN)). The literature on congestion control improvements for LFN is vast. It will not be covered further here, given the focus on AQM mechanisms in queues rather than congestion control algorithms in end-systems.\nReach Overload, Send ECN (ROSE)\nThe ROSE algorithm [WI99] is a queue marking variant of RED [RED] designed to improve its fairness while scaling to large networks. This is achieved by marking all packets whenever the queue size exceeds a threshold b and dynamically adapting b as follows:                if a packet would have been marked by SPSP, then decrease b, i.e. such that b←b−κε;        when a packet is marked then increase b, i.e. such that b←b+ε;where κ is a variable indicating how many marks are equivalent to a drop (1 or more) and ε is an infinitesimally small variable. This algorithm, which is based on actual queue length, can be shown to be fair to flows as it marks packets proportionally to the congestion they cause.        \nROSE is particularly relevant because it starts from the premise that an ECN mark need not be equivalent to a drop, instead the two being linearly related through the parameter κ. However, ROSE is a theoretical demonstration of a point, rather than a concrete algorithm proposal, so it does not actually specify the algorithm that should be used to drop packets rather than mark them, although perhaps it is implied that ROSE would be used with κ=1.\nMark and Drop Signals\nAs discussed in [open-ECN], the ECN specification for TCP/IP [RFC3168] stipulates that a packet should only be congestion marked if it would have been dropped, were it unmarkable. It is even stipulated that this assumption should be embedded in implementations, by stating that the ECT flag should only be checked after the decision has been made to drop a packet. Exactly mimicking drop behaviour is motivated by the need to provide incentives for hosts to switch to ECN capability when competing with unmarkable flows.\nThe [open-ECN] report was an openly published review of the ECN specification during its progress along the IETF's standards track. It agrees with the ECN specification that one ECN mark needs to be equivalent to one drop, but only within the best efforts service, for incremental deployment reasons. For other service disciplines than best efforts, it takes the position of Wischik [WI99] that one ECN mark need not be equivalent to one drop.\nData Centre TCP\nData-centre TCP (DCTCP), proposed by Alizadeh Attar et al. [ALI10] and also discussed in US patent application US2011/0211449 (Attar et al) has shown the considerable benefit of signalling congestion based on the instantaneous length of the queue rather than signalling congestion dependent on a moving average of the queue length. As a result, DCTCP keeps buffers nearly empty nearly all the time, resulting in very low and very predictable delay and extremely low loss-rates, both of which lead to significantly superior performance.\nDCTCP is an enhancement of TCP aimed at meeting the requirements of the applications running in data centres, so as to simultaneously achieve high throughput and burst tolerance as well as low latency. It relies on the usage of Explicit Congestion Notification [RFC3168]. The analysis of traffic patterns within Microsoft data centres have shown that even in the absence of the smoothing effects arising from multiplexing, DCTCP can achieve low queuing delay with low variance. DCTCP departs from earlier proposals by prescribing changes both at the senders and within network switches, which can be arranged to be deployed within a data centre where a single entity administers both:                Switch-side: the AQM in switches signals congestion by ECN marking packets based on instantaneous buffer length (i.e. no queue smoothing is carried out in switches as for example in RED). If the queue is greater than a threshold then it marks packets, otherwise it doesn't mark. This can be implemented with a RED AQM in commodity switches with a suitable choice of threshold parameters and the queue averaging parameter set to 1, which effectively turns off averaging.        Sender-side: it reacts to every congestion indication in proportion to congestion levels, i.e. the greater the number of ECN marks observed the greater the reduction in congestion window, Cwnd. The algorithm maintains average α of fraction F of packets marked in each RTT as follows:F=(Number of marked ACKs)/(Total ACKs)α←(1−g)α+g F (where 0<g<1 weight given to new samples)Cwnd←(1−α/2)Cwnd         \nDCTCP does not specify what the AQM should do for packets that do not support ECN, because ECN support is enabled by construction of the scheme throughout the data centre.\nPhantom Queue\nThe phantom queue [PR11] is aimed at achieving ultra-low latency solutions, particularly in data centres. The aim is to remove buffering delay from the data centre forwarding fabric by using a virtual queue on the switch (similar to PCN [PCN]) in combination with DCTCP [ALI10] on the end-systems.\nOther Techniques\nUS patent application US2003/0088690 (Zuckerman et al) relates to an AQM process which splits ECN and non-ECN packets into separate queues, then treats them differently. When congestion rises, the rate of ECN-marking on ECN-capable packets is increased. For non-ECN-capable packets, delay is introduced as an alternative to increasing loss.\nA paper entitled “A fair AQM scheme for aggregated ECN and non-ECN traffic” by Chong at al [ARQUA] describes a technique called ARQUA-DAB, which applies separate AQM algorithms to ECN and non-ECN traffic. ARQUA-DAB divides the capacity of the real link over two virtual buffers, each with service rates proportional to the number of ECN and non-ECN flows. It runs a substantially similar AQM algorithm in each virtual buffer, the only difference being the rates at which it drains the buffers.\nU.S. Pat. No. 7,139,281 relates to a method of active queue management for handling prioritised traffic in a packet transmission system.\nUS2011292801 relates to a proposed modification to the ECN protocol to allow a receiver terminal to exhibit some control of bandwidth share relative to other receiver terminals.\nUS2003112814 relates to a system and method for traffic management and regulation in a packet-based communication network which aims to facilitate proactive, discriminating congestion control on a per flow basis of packets traversing the Internet via use of a WRED algorithm that monitors the incoming packet queue and optimises enqueuing or discard of incoming packets to stabilise queue length and promote efficient packet processing.\nUS2004179479 relates to methods for alleviating traffic congestion by selectively discarding packets in a switch or router. In particular, it is directed to the problem of calculating average queue depth for RED in a manner that responds to rapid increases and decreases in the instantaneous queue depth."}
{"text": "The present invention relates to a method for the co-processing of waste rubber and carbonaceous material. More particularly, it relates to an improved method for liquefying waste rubber in mixture with a carbonaceous material such as coal to permit use of depolymerized rubber as a hydrogen donor solvent and carbon black reinforcing pigment as a liquefaction catalyst.\nIn the United States, worn and used rubber products are creating an increasingly large disposal problem. During the past few years, it is estimated that about 200 million used tires are discarded each year. Each tire contains about 15 pounds (7 kilograms) of organic materials and 2 kilograms of carbon black. If the rubber in these tires could be reprocessed with coal at about a 1-to-1 weight ratio, approximately 40,000 barrels per day of transportation fuels could be produced. In addition to the used tires, there is also a substantial amount of other waste rubber products that could be reprocessed. At the present time, most of these tires and rubber products are serving no useful purpose but are occupying significant space in waste dumps, vacant lots, and other vacant facilities.\nCoal liquefaction is typically achieved by heating the coal in the presence of a hydrogen donor solvent often with an added liquefaction catalyst. The most common solvent used are fractions derived from the liquified coal or heavy fractions from the petroleum industry in a co-processing operation. The liquefaction of coal is dependent on its hydrogenation which can be a function of the hydrogen donor solvent or the ability of the catalysts to promote reaction with hydrogen gas. Representative single and two-stage coal liquefaction processes are described generally in U.S. Pat. No. 4,874,506, and other issued U.S. Patents cited therein.\nTypically, particulate coal feed is slurried in a coal-derived recycle solvent and the resulting slurry is preheated to near the reaction temperature and fed with hydrogen into a catalytic reactor. The reactor operates at relatively high temperature and pressure conditions to hydroconvert and liquefy a major portion of the coal and thereby produce hydrocarbon gas and liquid fractions.\nUnfortunately, a large fraction of the coal liquefaction product is a residual oil containing preasphaltenes and asphaltene compounds having a low hydrogen content. It is these heavy, hydrogen-lean fractions that are often recycled for slurrying with the incoming coal to place an additional burden on the hydrogenation reactor."}
{"text": "The field of the disclosure relates generally to rotary machines, and more particularly, to a tip shroud of a turbine blade that includes notch features.\nAt least some known rotary machines include a compressor, a combustor coupled downstream from the compressor, a turbine coupled downstream from the combustor, and a rotor shaft rotatably coupled between the compressor and the turbine. Some known turbines include at least one rotor disk coupled to the rotor shaft, and a plurality of circumferentially-spaced turbine blades that extend outward from each rotor disk to define a stage of the turbine. Each turbine blade includes an airfoil that extends radially outward from a platform towards a turbine casing.\nAt least some known turbine blades include a shroud that extends from an outer tip end of the airfoil to reduce gas flow leakage between the airfoil and the turbine casing. An operational life cycle of at least some latter stage turbine blade tip shrouds may be limited by creep. Creep is the tendency of a material to deform over time when exposed to a combination of mechanical loading and high temperature. Turbine shroud creep rate may be greatly impacted by the high temperatures often seen at the shroud. To counter the effects of high temperatures, at least some known turbine blades include an internal cooling circuit, such as an interior tip shroud core cavity, or plenum, and/or passages that run transversely from the plenum toward the outer edges of the shroud. However, known tip shroud core plenums generally increase the complexity and expense of manufacture of the tip shroud, and impose design limits on other properties of the shroud, such as shape and thickness. Additionally, at least some known tip shrouds are shaped to reduce an area of the shroud that requires cooling. However, an ability to reduce the shroud area is limited by a need for the shrouds to cooperate structurally with adjacent shrouds."}
{"text": "1. Field of the Invention\nThe present invention relates to a plasma display panel, and more particularly, to a plasma display panel that improves the luminance efficiency by increasing a plasma density by forming a magnetic field within a discharge space.\n2. Description of the Related Art\nA Plasma Display Panel (PDP) display, which is a flat panel display, has excellent characteristics, namely, displays a high-quality image, is extremely thin and light, and provides a wide viewing angle, while having a large screen. In addition, a PDP display can be more simply manufactured than other flat panel displays and can be easily enlarged, such that the PDP display is spotlighted as a next-generation flat panel display.\nIn a 3-electrode surface discharge PDP, address electrode lines AR1, AG1, . . . , AGm, and ABm, front and rear dielectric layers, Y electrode lines Y1, . . . , and Yn, X electrode lines X1, . . , and Xn, phosphors, barrier ribs, and a MgO protective layer are disposed between front and rear glass substrates of the 3-electrode surface discharge PDP.\nThe address electrode lines AR1, AG1, . . . , AGm, and ABm are arranged in a predetermined pattern over the rear glass substrate. The rear dielectric layer is entirely coated over the address electrode lines AR1, AG1, . . . , AGm, and ABm. The barrier ribs are formed on the front surface of the rear dielectric layer to be parallel to the address electrode lines AR1, AG1, . . . , AGm, and ABm. The barrier ribs define discharge areas of each discharge cell and prevent optical crosstalk between adjacent discharge cells. The phosphors 16 are coated between barrier ribs.\nThe X electrode lines X1, . . . , and Xn and the Y electrode lines Y1, . . , and Yn are patterned on a rear surface of the front glass substrate so as to be orthogonal to the address electrode lines AR1, AG1, . . . , AGm, and ABm. The respective intersections define corresponding discharge cells. The X electrode lines X1, . . . , and Xn and the Y electrode lines Y1, . . . , and Yn are each comprised of a transparent electrode line of a conductive material, such as, Indium Tin Oxide (ITO), and a metal electrode line for increasing conductivity. For example, the X electrode line Xn is comprised of a transparent electrode line Xna and a metal electrode line Xnb, and the Y electrode line Yn is comprised of a transparent electrode line Yna and a metal electrode line Ynb. The front dielectric layer is entirely coated over the X electrode lines X1, . . . , and Xn and the Y electrode lines Y1, . . . , and Yn. The MgO protective layer for protecting the panel against strong electric fields is coated over the entire rear surface of the front dielectric layer. Discharge spaces are sealed with a gas therein for forming a plasma.\nIn the 3-electrode surface discharge PDP, not only the X electrode lines X1, . . . , and Xn, the Y electrode lines Y1, . . . , and Yn, but also the dielectric layer and the protective layer formed on the X and Y electrode lines exist on the front glass substrate. During discharge, visible rays emitted from the phosphors in the discharge spaces pass through the front substrate. However, the 3-electrode surface discharge PDP has a significant problem in that only about 60% of the visible rays are transmitted by the front substrate because of various components formed on the front substrate.\nAlso, in the 3-electrode surface discharge PDP, electrodes provoking the discharge are formed over the discharge spaces, namely, on the inner surface of the front substrate 10 through which the visible rays pass, such that the discharge is generated on the inner surface thereof and spreads. Hence, the 3-electrode surface discharge PDP has a low luminance efficiency.\nFurthermore, when the 3-electrode surface discharge PDP is used for a long period of time, charged particles of a discharge gas cause ion sputtering on the phosphors due to an electric field, thereby generating a permanent residual image."}
{"text": "1. Field of the Invention:\nThe present invention relates to a jumper wire inserting apparatus for inserting a short wire having substantially no resistance, namely, a jumper wire into a pair of holes of a printed circuit substrate.\n2. Description of the Prior Art\nHitherto, when a jumper wire inserting apparatus transforms a jumper wire responding to various distances, namely, pitches between two holes of a pair of holes of the printed circuit substrate and inserts the jumper wire into the holes by utilizing pushers and the like, both pushers are transferred in a horizontal direction responding to the various pitches of the holes. Therefore, the mechanism of the conventional jumper wire inserting apparatus is complicated and the accuracy of the transforming of the jumper wire is low.\nFurther, the horizontal direction transferring of the pushers and so on of the conventional jumper wire inserting apparatus does not have linkage with the operation of the jumper wire sending means for sending out a long wire for cutting to produce the jumper wire. Therefore, the accuracy of the operation of the conventional jumper wire inserting apparatus is low.\nAs above-mentioned, the conventional jumper wire inserting apparatus has various disadvantages."}
{"text": "In the related art, as an ignition system used in a combustion device of an internal combustion engine, a device, which includes a spark plug having a center electrode and a ground electrode, and forms a gap between both electrodes, a discharge power source such as an ignition coil supplying a voltage to the spark plug, and an ignition cable electrically connecting the spark plug and the discharge electrode, is known. In the above described ignition device, a high voltage is applied from the discharge power source to the spark plug via the ignition cable so that a spark discharge is generated in the gap of the spark plug. As a result, the fuel gas is ignited.\nIn recent years, a technique has been disclosed in which in order to improve an ignition ability, the spark discharge is generated by an AC power (a high frequency power) input from an AC power source to the gap, instead of the high voltage (for example, see JP-A-2009-8100)."}
{"text": "1. Technical Field\nEmbodiments of the present disclosure relates to circuit simulating systems and methods, and more particularly, to a computing device and a method for checking signal transmission lines of a circuit board.\n2. Description of Related Art\nA circuit board may be arranged with thousands of signal transmission lines. To ensure integrity of signals transmitted by the signal transmission lines, designs of the signal transmission lines should satisfy design standards. For example, a length of each line segment of a signal transmission line should satisfy a predetermined standard, so that each line segment of the signal transmission line has an appropriate impedance to ensure signal integrity. Spaces between neighboring signal transmission lines should also satisfy predetermined standards, to reduce crosstalk of the neighboring signal transmission lines. Therefore, it is necessary to incorporate design simulations and checks during the design and layout process of the circuit board. However, presently, obtaining information from a circuit board layout are often acquired manually. With the large quantity of signal transmission lines distributed on the circuit board, manual operation is not only time-consuming, but also error-prone."}
{"text": "1. Field of the Invention\nThe present invention relates to a carbazole derivative and an organic light emitting device using the same.\n2. Description of the Related Art\nAn organic light emitting device is a device in which a thin film including a fluorescent organic compound or a phosphorescent organic compound is interposed between an anode and a cathode. Further, electrons and holes are injected from the respective electrodes to generate exciton of the fluorescent compound or the phosphorescent compound, whereby light is emitted when the exciton return to a ground state is utilized.\nAccording to a study at Eastman Kodak company in 1987 (Appl. Phys. Lett. 51, 913 (1987)), there is reported a device having a function-separation type two-layer structure using ITO as an anode, a magnesium-silver alloy as a cathode, an aluminum quinolinol complex as an electron transporting material and a light emitting material, and a triphenylamine derivative as a hole transporting material. There has been reported a light emission of approximately 1,000 cd/m2 at an applied voltage of approximately 10 V.\nRecent progress of an organic light emitting device is remarkable, and the characteristics of the device enable a light emitting device with a high luminance at a low applied voltage, a variety of emission wavelengths, high-speed responsiveness, thin and light weight. From this fact, it is suggested that the device have potential to find use in a wide variety of applications.\nHowever, the present situation calls for optical output with even higher luminance or high conversion efficiency. In addition, many problems still remain to be solved regarding durability against the change over time due to long-term use, deterioration caused by atmospheric gas containing oxygen, moisture, or the like. Further, when considering application to a full color display, the present art is still insufficient against problems relating to the needs for light emission of blue, green, and red with a high color purity.\nTo solve the above-mentioned problems, there is proposed to use a carbazole derivative as a material for the organic light emitting device. As examples of the organic light emitting device using the carbazole derivative, there are given devices disclosed in “Journal of the American Chemical Society (2004), 126(19), 6035-6042,” and “Bulletin of the Korean Chemical Society (2004), 25(8), 1202-1206.”"}
{"text": "The present invention relates generally to vertebral implants and methods of use, and more particularly to implants and methods that include a flange that may be attached to a periphery of the implant body.\nThe spine is divided into four regions comprising the cervical, thoracic, lumbar, and sacrococcygeal regions. The cervical region includes the top seven vertebral members identified as C1-C7. The thoracic region includes the next twelve vertebral members identified as T1-T12. The lumbar region includes five vertebral members L1-L5. The sacrococcygeal region includes nine fused vertebral members that form the sacrum and the coccyx. The vertebral members of the spine are aligned in a curved configuration that includes a cervical curve, thoracic curve, and lumbosacral curve. Intervertebral discs are positioned between the vertebral members and permit flexion, extension, lateral bending, and rotation.\nVarious conditions may lead to damage of the intervertebral discs and/or the vertebral members. The damage may result from a variety of causes including a specific event such as trauma, a degenerative condition, a tumor, or infection. Damage to the intervertebral discs and vertebral members can lead to pain, neurological deficit, and/or loss of motion.\nVarious procedures include replacing the entirety or a section of a vertebral member, the entirety or a section of an intervertebral disc, or both. One or more replacement implants may be inserted to replace the damaged vertebral members and/or discs. The implants may further include bone growth material to facilitate fusion of the implant to one or both adjacent vertebral members. The implant should provide for housing the bone growth material, and prevent inadvertent removal of the material from the implant."}
{"text": "The present invention relates to fiber optic devices. In particular, it relates to a flexible microbend device that may be attached to an optical fiber and used in a sensor arrangement for measuring internal strain in a structure.\nOptical fibers are used to convey light between a light source and a light detector. Light in the fiber can be modulated by bending or otherwise distorting the fiber. This produces a modulated signal which can be picked up and processed by the light detector.\nMicrobend devices may be used to induce distortions in a fiber causing the light passing through the fiber to be modulated at a frequency corresponding to the force applied to the device. Typically, the device is applied to the fiber such that it surrounds the fiber.\nOne known device is a two jaw corrugated arrangement which squeezes an optical fiber to modulate the light signal passing therethrough. Macedo et al. (U.S. Pat. No. 4,342,907) describes such a device where an upper piece is firmly connected to an enclosure by a support. A lower piece was attached to an elastic membrane. When either a static pressure or dynamic time varying pressure reaches the membrane, it deflects pushing the lower piece, toward the upper piece, thus deforming the optical fiber which is held against the ridges of the device. They also describe another approach in U.S. Pat. No. 4,443,700 where two steel pieces having corrugated like surfaces with interleaving ridges form a vise which would be used to squeeze and thus deform the fiber. Such devices fail to have the compactness and flexibility desired for measuring forces internally in a structure.\nKramer (U.S. Pat. No. 5,193,129) introduced microbends by interweaving an optical fiber through rungs of a ladder-like structure, encapsulating it in a foam-like material, and surrounding it with a sheath. Light transmitted through the optical fiber is diminished to a value less than a threshold value upon the occurrence of microbending caused by pressure applied at any location along the length thereof. The rungs of the ladder structure are sized and spaced to provide a proper locus about which microbending may be produced. Sansone (U.S. Pat. No. 5,694,497) points out one of the deficiencies of this sensor by stating that it must use part of the structure it is embedded in to complete the sensor design. In fact, the utility of this type of sensor is limited in that it must be woven into the substrate or structure being measured and cannot be later repositioned as needed.\nUdd et al. (U.S. Pat. No. 5,118,931) introduced microbends into their sensor system by melting the optical fiber and pulling simultaneously to give the fiber a smaller diameter. Lengths of this fiber were spliced into the unmelted fiber to produce the sensor. Deformation of the fiber affects propagation of light therethrough, permitting detection based upon detected changes in light throughput. As with the Kramer device, the deficiency in this design is that once spliced into the fiber, there is no way to later reposition the structure.\nSansone (U.S. Pat. No. 5,694,497) describes an intrinsically self-deforming microbend deformer. In this device, the fiber is twisted about itself. At least one twisted section acts as an intrinsically self-deforming microbend deformer. The problem with this device is that there is no way to obtain a fixed period and/or thus be able to adjust the sensitivity of the sensor.\nAn object of the present invention is to provide a microbend device that may be removed and reattached to an optical fiber along its length.\nAnother object of the present invention is to provide a microbend device that is flexible.\nAnother object of the present invention is to provide a sensor that employs a flexible microbend device, such that the sensor is intensity-based.\nAnother object of the present invention is to provide a method for using a sensor that employs a flexible microbend device, such that the method is used to measure either radial or linear forces.\nBy the present invention, a flexible microbend device for attachment to an optical fiber is provided. The device comprises an upper bending element grid having at least one flexible element and a lower bending element grid having at least one flexible element. Each upper bending element grid alternately engages each lower bending element grid.\nWhen the device is attached to an optical fiber, it forms a fiber optic microbend sensor. The optical fiber is positioned between each flexible element of the upper bending element grid and each flexible element of the lower bending element grid such that each flexible element alternates between the upper and the lower bending element grids.\nIn use, the sensor is either attached to a surface of a structure or embedded in a host material. Optical power is introduced to the microbend sensor and the reflected light is monitored with a detector.\nAdditional objects and advantages of the invention will be set forth in part in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention will be obtained by means of instrumentalities in combinations particularly pointed out in the appended claims."}
{"text": "Activity in both humans and even animals depends largely on functioning, active muscles. As a consequence, sensing and monitoring of muscle activity is important for many reasons. For instance, in sports medicine, information about the fatigue level of a muscle is useful for preventing over-exertion of muscles, which might lead to muscle-based injuries. In addition, sensing muscle activity could be leveraged for precision training, allowing athletes to monitor progress on targeted muscle groups as needed. Finally, other ubiquitous applications such as muscle-controlled prosthetics could also be enabled by muscle activity data.\nThere are two main methods for detecting and monitoring muscle activity. The first is electromyography (EMG), which accurately measures the electrical potential of active muscle fibers. By amplifying and processing these electrical signals, information about muscle activation and activity can be inferred. The second method is mechanomyography (MMG). Mechanomyographic methods use vibration transducers to detect skeletal muscle activity. A brief discussion of the physiology of skeletal muscle vibrations that enable mechanomyography based systems is provided herein.\nFIG. 1 illustrates an example skeletal muscle 100 and in particular the quadriceps and the tendons 106a, 106b and muscle belly 102 (including muscle fibers 104) are shown. Skeletal muscles 100 generally consist of three parts: a muscle belly 102, composed of muscle fibers 104, and two tendon ends 106a, 106b. Activating skeletal muscles during physical activity causes the muscle fibers 104 to contract, resulting in mechanical vibrations due to: tremors of the human motor system, muscle fibers 104 sliding against each other and artifacts detected when a muscle belly's 102 circumference increases during muscle contractions and decreases during relaxation. These skeletal muscle vibrations are carried through to and detected at the skin level. The vibrations vary depending on the muscle fatigue status, muscle group size, and location in the body e.g., skeletal muscle proximity to bones of different mass and size.\nAs a method for inferring muscle activity, EMG while accurate, has the major setback that it is difficult, invasive and even painful to use. It requires perfect contact with the skin, necessitating the use of gels in the case of surface EMG, or the insertion of small needles into the muscle, in the case of fine-wire EMG. Mechanomyography methods on the other hand, work by externally detecting the mechanical vibrations occurring in the muscle belly of skeletal muscles when the muscle is activated using vibration transducers. As such they are less invasive and more convenient for muscle monitoring use in a dynamic environment. Current vibration transducer technologies include fabric stretch sensors, force sensitive resistors (FSR) and a combination of any of or all of available inertial sensors including accelerometers, gyroscopes and magnetometers."}
{"text": "As described in Patent Document 1, a tablet cassette has been structured as follows: a rotatably-driven rotor is disposed in a bottom of a cassette body; the outer peripheral surface of the rotor is formed with a plurality of guide grooves in an axial direction and a circumferential groove in a circumferential direction; the guide groove is communicated with a tablet outlet formed in the bottom of the cassette body; a partitioning member is attached to the outer surface of the cassette body so as to enter the circumferential groove of the rotor; and the lowermost tablet and another tablet placed thereon, which are introduced to the guide groove, are separated from each other by the partitioning member so that only the lowermost tablet is dispensed from the tablet outlet. Patent Document 2 further describes a similar partitioning member that is inserted in an insertion hole of a cassette body and is secured by a cover. The aforementioned tablet cassettes need to use the rotor that has a different guide groove depending on the size of the tablet, and also need to change the attachment position of the partitioning member.\nThus, various proposals have been made so as to dispense a plurality of types of tablets having different sizes using the same tablet cassette.\nPatent Document 3 discloses a tablet cassette configured as follows: a rotor includes a core body and an outer cylindrical body releasably engaging the core body; the outer cylindrical body having a guide groove with different width and depth is replaceable depending on the size of the tablet to be dispensed; a partitioning member is configured to be movable up and down by tightening or loosening a bolt; and the position of the partitioning member is changed depending on the size of the tablet to be dispensed. However, the above-described structure requires two complicated tasks, i.e., a change of the outer cylinder bodies and replacement of the partitioning member. Patent Document 3 also discloses that a ring body is fitted between the outer surface of the rotor and the inner surface of the cassette body to reduce the cross-section of the guide groove, and that a groove-depth adjuster movable in a radial direction is provided on a bottom of the guide groove so as to adjust the depth of the guide groove. However, it is difficult to mount and fix the ring body, and moving and fixing the groove-depth adjuster for each guide groove are complicated and time-consuming. Since the outer cylindrical bodies and the ring bodies are changed according to the shape of the tablet, replacement parts therefor are required.\nPatent Document 4 discloses the following structure: a plurality of insertion holes are provided in a cassette body, and a partitioning member is fitted and fixed to the insertion hole located at a height corresponding to the size of a tablet; alternatively, a plurality of insertion holes and positioning holes are provided, the partitioning member is inserted in the insertion hole located at a height corresponding to the size of a tablet, and a pin provided in the partitioning member is fitted into the positioning hole to fix the partitioning member. However, the above-described structure has a problem that it cannot adjust the entry amount of the partitioning member to a guide groove.\nIn Patent Documents 1 to 4, the lowermost tablet and another tablet placed thereon, which are introduced to the guide groove, are separated from each other by the partitioning member, and these tablets make a sliding contact with the partitioning member as the rotor rotates. This results in a problem that the tablets are easily damaged."}
{"text": "This invention relates to thermoformed thermoplastic truck bedliners in general and to bedliners with structure for engaging load restraining members in particular.\nAlthough long used in agricultural and commercial applications, pick-up trucks with open cargo beds have also become increasingly popular as personal and family vehicles. The commercial vehicle owner, although concerned with performance and cost effectiveness of the vehicle, often considers the vehicle as a traveling company advertisement or symbol. A neat and well-maintained vehicle is more likely to favorably impress customers. Owners of personal vehicles, while concerned with functionality and efficiency, are also concerned with pride of ownership, personal image, and outward appearances. In addition, the automotive enthusiast desires to maintain his vehicle in as close to a \"like new\" condition as possible.\nTruck bedliners of plastic or rubber are commonly employed to protect the pickup truck cargo bed. These truck bedliners are available in a wide variety of configurations to suit the wide variety of available trucks. One-piece thermoformed thermoplastic truck bedliners provide a cost effective means of protecting the truck cargo bed.\nTo facilitate the handling of cargo, many truck cargo boxes have features which stabilize and control shifting of cargo in the truck box. For instance, the box may have stake holes along the upper peripheral edges so that boards may be inserted into the sides of the truck box, so extending the sides to retain light-weight, high-volume loads. Some truck boxes have supports on the side walls for horizontal boards, so that wide loads may be supported above the interiorly extending truck wheel wells on horizontal boards.\nSome trucks have features formed into the sides of the truck to support a plank or \"2.times.4\" board so it extends upright across the width of the bed. Boards extending across the floor of the bed of the cargo box are useful to prevent the shifting of cargo as the truck accelerates and brakes.\nTruck bedliners have been developed which extend existing features in truck cargo boxes and which provide features such as horizontal board pockets for use in truck boxes that lack them. For instance, one type of truck bedliner provides a series of board slots formed by vertical ribs which extend into the interior of the cargo box. These series of board slots allow a vertical 2.times.4 board to be adjustably positioned among the slots to secure the load into place.\nBedliners come in two general types, so-called \"over-the-rail\" liners which extend up and cover the upper edges of the cargo box, and so-called \"under-the-rail\" liners in which side-walls of the liner terminate under the upper flanges or rails of the cargo box.\nBoard holder pockets have been used for a number of years to support boards horizontally to prevent shifting of cargo. Board holder pockets, however, open upwardly, and hence retain the possibility that a board held therein may be rotated or shifted out of the pocket if subjected to an upwardly-directed force.\nA thermoformed truck bedliner is needed which will accept commonly sized timbers to extend across the truck bed and which will support the timbers against loads applied by shifting cargo."}
{"text": "Dietary supplements containing viable probiotic bacteria are increasing in popularity in the marketplace as their health benefits become recognized. Reported benefits range from alleviating constipation and diarrhea to reducing various intestinal infections such as those caused by rotaviruses, pathogenic E. coli and Helicobacter pylori as discussed in U.S. Pat. Nos. 7,090,840 and 7,029,669. Beneficial species of Lactobacillus and bifidobacteria are among the widely recognized probiotics, and strains of these bacteria that are capable of colonizing the intestinal tract are advantageous (see U.S. Pat. Nos. 7,150,986 and 6,887,465).\nColonization may involve physical attachment to epithelial cell surfaces such as those of the microvilli in the ileum section of the small intestine, or simple domination of the contents of the cecum, or adherence within the mucin layers of the colon. There is much to learn about how probiotics colonize the mammalian intestinal tract but it has been established that when colonization occurs, more probiotic microorganisms appear in the feces and this correlates with more probiotic microorganisms in both the proximal and distal sections of the intestinal tract. See Muralidhara et al, 1997, Journal of Food Protection Vol. 40, No. 5, Pages 288-295. Young pigs were used to demonstrate the relationship between colonization and fecal probiotic counts in Muralidhara's work; the porcine intestinal tract is very similar in physiology to the human intestinal tract and similar studies with young pigs are presented in examples of the present invention to demonstrate its effectiveness.\nColonization that results in the competitive exclusion of pathogenic microorganisms is particularly beneficial and can occur when probiotics occupy most of the intestinal attachment sites and are encouraged to produce lactic acid and other antimicrobial compounds. Effective intestinal colonization by probiotics depends on the availability of proper microbial nutrition that must be provided by the diet. See Gibson et al 1995, Gastroenterology 106: 975-982; Christl et al, 1992, Gut 33: 1234-1238 and Reid et al, 1990, Clin. Microbiol. Rev. 3: 335-344. However, normal diets do not provide nutrients that necessarily benefit probiotics so there is a need to fortify the diet with such nutrients. Prebiotics are one class of microbial nutrients that are currently popular in the marketplace. They are typically certain oligosaccharides that are not digested in the small intestine but serve as nutrients for select probiotic bacterial genera, e.g., bifidobacteria, when they arrive in the colon. An article titled “Prebiotics Enhance Gut Health” by Laura Brandt is available online at the nutrasolutions website; another titled “Prebiotics: A More Reliable Way to Increase Gut-Friendly Bacteria” by Dr. James Meschino is available at the chiroweb website, in the archival section. Today, a variety of functional foods are fortified with prebiotics such as fructooligosaccharides (FOS) and inulin, in an effort to provide probiotic stimulation in vivo. This practice adds cost to the foods being fortified and is not very effective in stimulating Lactobacillus probiotics. Thus, there is a need for reducing the cost of prebiotics while providing for both Lactobacillus and Bifidobacterium stimulation.\nAlthough enzymes have been used to generate prebiotics under laboratory conditions followed by subsequent feeding of the preformed prebiotics to achieve probiotic stimulation (see U.S. Pat. Nos. 6,791,015 and 6,730,502), no one has suggested using enzymes to generate these effects in vivo. U.S. Pat. No. 5,817,350 discloses the use of the prebiotic enzymes cellulose, amylase and hemicellulase, for use as dietary supplements, but not use of these enzymes to stimulate administered probiotics, or enhancement of their prebiotic effect by addition of polysorbate compounds. U.S. patent application 20010031276 discloses the use of polysorbate compounds as feed additives for ruminant animals but does not relate to their use in combination with prebiotics and probiotics, or their use in non-ruminant animals.\nIf polysorbate surfactants are combined with enzymes, water activity (Aw) becomes important. High levels of Aw, e.g., Aw>0.04, can significantly destabilize the shelf life of certain enzymes. Polysorbates are viscous, sticky, non-aqueous liquids that are not available in dry form. Before incorporation into enzyme formulations, polysorbates must be rendered into dry powders and the Aw of any formulation containing them should be below 0.04.\n“Internal Probiotic Culture” or IPC, as used herein, is the totality of viable probiotic microorganisms present in the intestinal tract at any given time: i.e., the sum of probiotic microorganisms adhering to intestinal epithelial cells, mucous and mucin layers, ingested food and waste material. It is desirable to enhance the IPC of humans and mammals. IPC can be estimated from the total viable Lactobacillus and/or Bifidobacteria counts (colony forming units) present in fresh fecal matter."}
{"text": "1. Field of the Invention\nThis invention relates to the fabrication of a thin-film magnetic head. In particular, it relates to a method of fabricating a stitched writer portion of such a head that is suitable for high data-rate recording of magnetic information.\n2. Discussion of the Related Art\nThin film magnetic heads are used to record and retrieve data stored in the form of small magnetized regions on disks and tapes. These heads contain a read portion, which is typically a shielded magnetic field sensor of a giant magneto-resistive (GMR) type, and a write portion, consisting of a magnetic pole and yoke structure inductively energized by current carrying coils.\nThe need to record and retrieve data stored with increasing area densities and at increasingly higher data rates has necessitated the design of very narrow write heads with high linear resolution. The narrow write head allows data to be stored in correspondingly narrow tracks. The high linear resolution, which is achieved by the formation of a thin write gap and a strong magnetic field, allows the storage of more data per unit length of track.\nAn important advance in the fabrication of write heads that help to meet these stringent requirements is the so called xe2x80x9cstitched write head design.xe2x80x9d This design embodies a simplified method of forming a narrow upper magnetic pole piece by joining or xe2x80x9cstitchingxe2x80x9d together two separately deposited pole pieces, the upper pole tip and the upper pole yoke, along a pedestal formed from a portion of the upper pole tip photoresist mask. In contrast to methods that form the upper pole piece and yoke monolithically, the stitching process allows the narrow tip section to be formed within a thinner photoresist mask, which is highly advantageous in the context of the fabrication process.\nThe performance of the stitched writer design and the advantages inherent in the nature of its fabrication process can be further improved by overcoming several deficiencies which limit its applicability to high data-rate recording. In particular, the stitched writer head as formed by the methods of the present art has a lengthy throat region along which the magnetic flux is diminished. Thus, in order to maintain an acceptable flux level for recording at the air-bearing surface (ABS), the design requires a high saturation writing current. This, in turn, causes side erasures and adversely affects data already written on adjacent disk areas.\nSeveral methods have been advanced to improve the flux characteristics of the stitched writer design and, indeed, to improve such other designs as have been suggested to meet the need for high recording resolution. Chang et al. (U.S. Pat. No. 5,805,391) teach a method for forming a write head with a recessed stitched yoke on a planar portion of an insulating layer. The contour of said insulating layer minimizes the flux leakage between the yoke portion and the pole tip portion stitched to it. Two inventions of Chen et al. (U.S. Pat. No. 5,652,687 and U.S. Pat. No. 5,802,700) teach methods for forming a write head having a U-shaped notch which opens out at the air-bearing surface and between whose ends is formed the upper pole tip. The structure so formed provides a parallel path for conducting the magnetic flux to said pole tip, thereby allowing the proper flux concentration for writing in narrow trackwidths. The invention of Cole et al. (U.S. Pat. No. 5,452,164) teaches a method of forming a pole tip structure that lays both above and below the write gap and has a track width of less than 1 micron. The symmetry of this arrangement supports good flux transfer with no leakage. The invention of Santini (U.S. Pat. No. 5,621,596) provides a method for improving the resolution of the photolithography process that is used to define the shape of the pole tip. Specifically, the invention teaches a method of reducing light reflections in the photolithography process by moving the point at which the pole yoke flares sufficiently far back from the position of the pole tip so that reflections from the flare do not interfere with the definition of the pole tip.\nAn alternative approach to improving the flux characteristics of the stitched pole design involves concentrating the flux by improving the geometry of the junction between the upper pole tip and upper pole yoke. This is presently being achieved by a design (FIG. 1) that incorporates a recessed upper pole yoke (FIG. 1 (10)) and a xe2x80x9cstepxe2x80x9d (FIG. 1 (12)) formed of non-magnetic material positioned in the rear underside of the pole tip where it contacts the write gap layer. The step effectively channels the flux through a smaller area, thereby increasing its intensity across the write gap region (FIG. 1 (14)) and enhancing the writing process. Although the use of the step provides a significant improvement in the stitched head design, its placement in the pole tip will ultimately prove to be disadvantageous as thinner pole tips and shorter throat heights are increasingly necessitated for high data-rate applications. The present invention teaches a method for fabricating a recessed yoke, stitched head writer with the flux concentrating advantages of the step formation, but without the disadvantages of the step placement in the upper pole tip.\nA first object of this invention is to fabricate a stitched pole magnetic write head that is capable of recording magnetic data at high rates and increased densities.\nA second object of this invention is to fabricate a stitched pole magnetic write head that has a lower saturation write current.\nA third object of this invention is to fabricate a stitched pole magnetic write head with improved nonlinear transition shift performance.\nA fourth object of this invention is to fabricate a stitched pole magnetic write head having improved overwrite performance.\nA fifth object of this invention is to fabricate a stitched pole magnetic write head that significantly reduces the problem of side erasure.\nA sixth object of this invention is to fabricate a stitched pole magnetic write head with a reduced effective throat height and upper pole thickness.\nThese objects will be achieved by means of a novel modification of the present stitched pole, stepped pole-piece fabrication process. The proposed modification is the removal of the step from the underside of the upper pole piece and its effective repositioning in the form of a xe2x80x9crecessed step,xe2x80x9d a non-magnetic spacer embedded into an upper layer of the lower pole of the write head, said spacer thereby being placed beneath the write gap layer of said write head. The flux concentrating region so formed is less sensitive to the position of the upper pole piece, thereby allowing the formation of a shorter throat height as will be required in future high data rate applications."}
{"text": "Various electronic parts mounting apparatus are known. For example, as disclosed in Japanese Utility Model Application Laid-Open No. 57-48672, previously filed by the present applicant, a tipped part seal-in tape is wound and accommodated in a tape reel, the tape being intermittently drawn and delivered to a fixed part-removing station. The tipped parts are suctioned off the tape by a plurality of reciprocally movable suction nozzles which are provided below a lower portion of the peripheral edge of an intermittently rotating turntable. The tipped parts are mounted on a table which is movable and controllable in an X-Y direction for positioning and identifying the tipped parts suctioned and held by the suction nozzles.\nPositioning apparatus for tipped parts held by the suction nozzles are described, for example, in Japanese Patent Application Laid-Open No. 56-24,997 and Japanese Utility Model Application Laid-Open No. 58-136,288 Specifications, in which the tipped parts are positioned according to an arranging mode on a print substrate. For purposes of this application, \"print substrate\" means a printed circuit board. The tipped parts are urged from various directions by positioning pawl members disposed in a fixed positioning station. Japanese Patent Application Laid-Open No. 61-168,298 Specification, discloses an apparatus in which the position and shape of the tipped parts are recognized in advance by optical means. The suction nozzles are then rotated to properly position the tipped parts and simultaneously correct deviations in the positioning of the rotating nozzles at the time of suction.\nHowever, in the conventional positioning apparatus and in the apparatus disclosed in Japanese Patent Application Laid-Open No. 56-24,997 Specification, the positioning and rotation of the nozzles to adjust the direction of the tipped parts cannot be effected at the same time, which limits the speed of the apparatus. An apparatus disclosed in Japanese Utility Model Application Laid-Open No. 58-136,288 shows in FIG. 27 a V-shaped notched groove 201 which is formed at the end of a suction nozzle 200 to suction and hold a cylindrical tipped part. When a center between the V-shaped notched groove 201 at the end of the suction nozzle 200 and each of the positioning pawl members 203 and 204 is deviated from each other, the cylindrical tipped part 202 tends to be positioned by the positioning pawl members 203 and 204 in an unstable state where the tipped part 202 does not properly fit in the V-shaped notched groove 201. In addition, when the tipped part 202 is released from each of the positioning pawl members 203 and 204 after completion of the positioning, the tipped part 202 may become disoriented due to movement in such a way that the tipped part 202 is fitted into the notched groove 201 under the influence of the suction force of the suction nozzle 200. For this reason, it is difficult to adjust the direction of the notched groove 201 at the end of the suction nozzle and the direction of each of the positioning pawl members 203 and 204 for providing precise coincidence, which has to be carried out for each suction nozzle. Moveover, the apparatus must be replaced according to the shape and size of the tipped parts.\nIn the apparatus disclosed in Japanese Patent Application Laid-Open No. 61-168,298 Specification, when the rotational center between the suction nozzles is deviated, the position of the tipped parts after rotation becomes deviates, and therefore, a troublesome precise adjustment of the rotational center of each of the suction nozzles needs to be carried out."}
{"text": "1. Field of the Invention\nThe present invention relates to a thermal ink-jet ink having excellent storage stability and an ink cartridge for containing the ink.\n2. Description of the Related Art\nAn ink-jet recording method is a recording method involving causing a small ink droplet to fly to plain paper or a dedicated glossy medium to form an image. The method has become rapidly widespread in association with advance of a reduction in price of a recording apparatus and an increase in recording rate of the apparatus. In particular, there has been a growing need for photograph picture quality because a digital camera has become widespread, so an additional improvement in image quality and printing at an additionally high speed have been requested. Accordingly, more sophisticated techniques than the conventional ones have been requested, and examples of the techniques include: a reduction in size of an ink droplet; an increase in density at which nozzles are arranged; an increase in length of a head in association with an increase in number of nozzles; and control of ejection of an ink droplet.\nOn the other hand, a thermal ink-jet recording mode is a mode involving: foaming ink by utilizing thermal energy; and ejecting the ink to a recording medium. The mode enables high-speed, high-density, high-definition, high-quality recording, and is suitable for colorization and a reduction in size of a recording apparatus. A general head to be used in the recording mode includes: a substrate for an ink-jet recording head on which a heat element for foaming ink and wiring to be electrically connected to the heat element are produced; and a flow path for ejecting the ink on the substrate.\nIn addition, the substrate for an ink-jet recording head is modified in various ways in order that electrical energy to be inputted may be saved and a reduction in lifetime of the substrate resulting from destruction of a heat generating portion in association with the foaming of ink may be prevented. In particular, a protective layer for protecting a heat element placed between a pair of wiring patterns from ink is modified in many ways.\nIt is advantageous for the protective layer to have a high thermal conductivity or a small thickness from the viewpoint of thermal efficiency. However, it is advantageous for the protective layer to have a large thickness from the viewpoint of protection of wiring to be connected to the heat element from ink. Accordingly, the thickness of the protective layer must be set to an optimum thickness from the viewpoints of energy efficiency and reliability. In particular, a layer in contact with ink are affected by both of cavitation damage due to the foaming of the ink, that is, mechanical damage, and damage due to a chemical reaction with an ink component at high temperature, that is, chemical damage, so influences of the mechanical damage and the chemical damage need to be sufficiently taken into consideration.\nIn view of the foregoing, the protective layer of an ink-jet substrate generally has a layer having high stability against mechanical damage and chemical damage as an upper layer (i.e., layer in contact with ink), and an insulating layer for protecting wiring as a lower layer. To be specific, in general, a Ta layer, which is a layer extremely stable both mechanically and chemically, is formed as the upper layer, and an SiN layer, SiO layer, or SiC layer, which can be easily formed by using an existing semiconductor producing device and is stable, is formed as the lower layer.\nDetailed description of the foregoing is as described below. An SiN layer having a thickness in the range of from about 0.2 μm or more to 1 μm or less is formed as a protective layer on wiring, and then a protective layer as an upper layer is formed. A Ta layer has a thickness in the range of from 0.2 μm to 0.5 μm, which is referred to as a cavitation resisting layer, is formed as the upper layer because the layer serves as a layer to cope with cavitation damage. With this constitution, compatibility between the lifetime and reliability of a heat element of an ink-jet substrate is achieved.\nJapanese Patent Application Laid-Open No. H05-330048 discloses, as an ink-jet technique using a thermal head, an ink-jet head using a material containing Si, N, or Ir at a specific ratio for a heat element to improve durability and thermal conversion efficiency. There is also disclosed an ink containing a chelating reagent at a specific concentration from the viewpoints of improvement in ejection durability of a thermal head and in suppression of a kogation to be deposited at a heat generating portion due to ejection (see Japanese Patent Application Laid-Open No. H06-093218). There is also disclosed an ink containing an ammonium salt of an acid having a methyl group or a methylene group, and a carboxyl group (see Japanese Patent Application Laid-Open No. 2002-012803). Those conventional arts each relate to suppression of erosion of a Ta layer in association with ejection duration or the suppression of a kogation to be deposited on the Ta layer when the Ta layer is arranged as the surface layer of the protective layer of a heat generating portion. In addition, ink contains a specific compound at a specific concentration, whereby a balance between the deposition of a kogation and the erosion of the Ta layer due to ejection duration is optimized for the lengthening of the lifetime of the ink."}
{"text": "The present invention relates to a new 5-enolpyruvylshikimate-3-phosphate synthase (or EPSPS) which displays increased tolerance with respect to herbicides which are competitive inhibitors with respect to phosphoenolpyruvate (PEP) of EPSPS activity. This more tolerant EPSP synthase possesses at least one xe2x80x9cthreonine by isoleucinexe2x80x9d substitution. The invention also relates to a gene coding for such a protein, to plant cells transformed by chimeric gene constructions containing this gene, to the plants regenerated from these cells and also to the plants originating from crossing using these transformed plants.\nGlyphosate, sulfosate and fosametine are broad-spectrum systemic herbicides of the phosphonomethylglycine family. They act essentially as competitive inhibitors of 5-enolpyruvylshikimate-3-phosphate synthase (EC 2.5.1.19) or EPSPS with respect to the PEP (phosphoenolpyruvate). After their application to the plant, they are translocated in the plant where they accumulate in the rapidly growing parts, in particular the cauline and root spices, causing damage to the point of destruction of sensitive plants.\nPlastid EPSPS, the main target of these products, is an enzyme of the pathway of biosynthesis of aromatic amino acids, which is encoded by one or more nuclear genes and synthesized in the form of a cytoplasmic precursor, then imported into the plastids where it accumulates in its mature form.\nThe tolerance of plants to glyphosate and to products of the family is obtained by stable introduction into their genome of an EPSPS gene, of plant or bacterial origin, which is mutated or otherwise in respect of the characteristics of inhibition by glyphosate of the product of this gene. In view of the mode of action of glyphosate and the degree of tolerance to glyphosate of the product of the genes which are used, it is advantageous to be able to express the product of the translation of this gene so as enable it to be accumulated in substantial amounts in the plastids.\nIt is known, for example from U.S. Pat. No. 4,535,060, to confer on a plant a tolerance to a herbicide of the above type, especially N-phosphonomethylglycine or glyphosate, by introducing into the genome of plants a gene coding for an EPSPS carrying at least one mutation that makes this enzyme more resistant to its competitive inhibitor (glyphosate) after localization of the enzyme in the plastid compartment. These techniques, however, need to be improved in order to obtain greater reliability in the use of these plants under agricultural conditions.\nIn the present description, xe2x80x9cplantxe2x80x9d is understood to mean any differentiated multicellular organism capable of photosynthesis, and xe2x80x9cplant cellxe2x80x9d is understood to mean any cell originating from a plant and capable of constituting undifferentiated tissues such as calluses or differentiated tissues such as embryos or plant parts or seeds.\nThe subject of the present invention is the production of transformed plants having increased tolerance to herbicides of the phosphonomethylglycine family, by regeneration of cells transformed by means of new chimeric genes containing a gene for tolerance to these herbicides.\nThe subject of the invention is also a chimeric gene for conferring on plants increased tolerance with respect to a herbicide having EPSPS as its target, comprising, in the direction of transcription: a promoter region, optionally a transit peptide region, a sequence of a gene coding for a glyphosate tolerance enzyme and an untranslated polyadenylation signal region at the 3xe2x80x2 end, characterized in that the glyphosate tolerance gene contains, relative to the gene from which it is derived, a xe2x80x9cthreonine 102 by isoleucinexe2x80x9d substitution in the xe2x80x9caroAxe2x80x9d (EPSPS) region. Preferably, it comprises, in addition, in the same region, a xe2x80x9cproline 106 by serinexe2x80x9d substitution. These substitutions can be introduced or be present in an EPSPS sequence of any origin, in particular of plant, bacterial, algal or fungal origin.\nThe transit peptides which can be used in the transit peptide region can be, known per se, of plant origin, for example originating from maize, sunflower, pea, tobacco or the like. The first and the second transit peptide can be identical, similar or different. They can, in addition, each comprise one or more transit peptide units according to European Patent Application AP 0 508 909. It is the role of this characteristic region to permit the release of a mature and native protein, and especially the above mutated EPSPS, with maximum efficacy in the plasmid compartment.\nThe promoter region of the chimeric gene according to the invention may be advantageously composed of at least one gene promoter or promoter fragment which is expressed naturally in plants (tubulin, introns, actin, histone).\nThe untranslated transcription termination signal region at the 3xe2x80x2 end of the chimeric gene may be of any origin, for example of bacterial origin, such as that of the nopaline synthase gene, or of plant origin, such as that of the Arabidopsis thaliana histone H4A748 gene according to the European Patent Application (European Application 633 317).\nThe chimeric gene according to the invention can comprise, in addition to the essential portions above, at least one untranslated intermediate (linker) region, which can be located between the different transcribed regions described above. This intermediate region can be of any origin, for example of bacterial, viral or plant origin.\nIsolation of a cDNA coding for a maize EPSPS:\nThe different steps which led to the obtaining of maize EPSPS cDNA, which served as substrate for the introduction of the two mutations, are described below. All the operations described below are given by way of example, and correspond to a choice made from among the different methods available for arriving at the same result. This choice has no effect on the quality of the result, and consequently any suitable method may be used by a person skilled in the art to arrive at the same result. Most of the methods of engineering of DNA fragments are described in xe2x80x9cCurrent Protocols in Molecular Biologyxe2x80x9d Volumes 1 and 2, Ausubel F. M. et al., published by Greene Publishing Associates and Wiley-Interscience (1989) (hereinafter, references to protocols described in this work will be designated xe2x80x9cref. CPMBxe2x80x9d). The operations relating to DNA which were performed according to the protocols described in this work are especially the following: ligation of DNA fragments, treatment with Klenow DNA polymerase and T4 DNA polymerase, preparation of plasmid and of bacteriophage xcex DNA, either as a minipreparation or as a maxipreparation, and DNA and RNA analyses according to the Southern and Northern techniques, respectively. Other methods described in this work were followed, and only significant modifications or additions to these protocols have been described below."}
{"text": "This invention concerns so called “knee bolsters”, which are structures installed in automotive vehicles forward of the driver's seat positioned to be engaged by the driver's knees if the person submarines below the steering wheel mounted air bag during a collision.\nKnee bolsters sometimes employ crushable or deformable impact absorbing structures which restrain the driver by contact by his or her knees, the crushing of the structure absorbing the energy of the person's momentum over a range of deformation to lower the stress on the person's knees and the load transmitted to the thigh bone.\nDeformable structure knee bolsters are designed to be crushed by the person as the person is decelerated by the collision in order to reduce or prevent injury to the person.\nAnother requirement is the prevention of damage to or interference with the steering column mechanism during crushing of the knee bolster impact absorbing structure so as to allow the steering column to steer axially and absorb the upper body impact and to allow steering control to be maintained after a collision, if possible.\nCosts are always a problem in automotive design, and the knee bolsters have sometimes been assembled from a number of components, representing an assembly cost.\nLarge aluminum sections have been proposed to be extruded to eliminate assembly costs, with portions cut out from the large extrusions. Large extrusions are themselves costly such that these structures are still somewhat expensive to manufacture.\nOther knee bolsters have involved crush cells which have curved webs which are deformed, but such webs have a geometry which results in a buckling collapse of the web, greatly reducing the structure's ability to absorb energy. This can result in a great reduction in the restraining resistance of the structure. Such collapse can also produce gross bulging which can interfere with proper operation of the steering mechanism. Also, such a collapse mode can cause the surface impacted by the knees to directly hit the steering column and hinder the proper functioning of the steering column.\nIt is the object of the present invention to provide a deformable structure type knee bolster which is low in cost yet reliably providing adequate resistance to deformation over a range of crushing movement and prevents the development of any interference with the steering column mechanism."}
{"text": "Meat skinning machines conventionally have a toothed gripping roll located adjacent a sharpened skinning blade for separating the skin from meat, fish and poultry products. The products to be skinned are moved manually into the gripping roll and towards the blade as the skinning operation takes place. Obviously, any contact between the operator's hands and the gripping roll and blade is likely to invite injury to the operator.\nSafety circuits have been incorporated into skinning machines whereby the power to the operating components of the machine will immediately stop whenever the hands of the operator move into and contact a danger zone in the immediate area of a gripping roll and blade. Typical safety circuits are disclosed in my copending applications Ser. No. 426,668 filed Oct. 26, 1989 and Ser. No. 447,156 filed Dec. 7, 1989. The safety circuits of these two applications are incorporated herein by reference.\nOne of the difficulties in providing a safety circuit and apparatus for meat skinning machines dwells in the matter of instantaneously stopping the power mechanism when the operator inadvertently contacts components in the danger zone of the machine. The inertia of electrical motors contributes to this problem of obtaining instantaneous stoppage of the machine when a dangerous situation arises.\nIt is therefore a principal object of this invention to provide a safety means for meat skinning machines which will substantially instantaneously stop the gripping roll of the machine when a safety circuit is actuated without also having to deactivate the power source of the machine.\nA further object of this invention is to provide a safety means for meat skinning machines which is adaptable for use with a plurality of different safety circuits.\nA still further object of this invention is to provide a safety means for meat skinning machines which is reliable and easily serviced.\nThese and other objects will be apparent to those skilled in the art."}
{"text": "This invention relates to apparatus for continuously variable rotary power transmission. More particularly, this invention relates to such apparatus which incorporate variously sized and shaped counter rotating friction wheels.\nA known continuously variable power transmission comprises a conical wheel which is frictionally paired with a second wheel, the axes of rotation of the conical and second wheels typically being aligned so that they lie within a single flat reference plane. In such configuration, while the points of frictional contact between the conical and second wheels are near the conical wheel\"\"s base, rotation of the conical wheel counter rotates the second wheel at a high speed. Conversely, frictional contact near the conical\"\"s wheel\"\"s point results in counter rotation of the second wheel at a relatively low speed. Such an assembly provides continuously variable rotary power transmission as a result of a potentially infinite number of positions of the second wheel between the conical wheel\"\"s point and base.\nIn order for the above exemplary conical wheel transmission to effectively transmit rotary power, an acceptable level of friction must exist at the points of contact between the two wheels. A known means for enhancing wheel driving friction at such contact point forms one or both of the wheels from a material having a high coefficient friction, such as rubber. However, high friction materials, such as rubber, are often undesirably utilized because they lack durability. Another known means of enhancing wheel driving friction at the contact point is to widen the contact face of the second wheel, resulting in a lengthening of its zone of frictional contact with the conical wheel. However, in such configuration, opposing sliding friction force vectors invariably arise on opposite sides of a centrally located useful static friction zone. The resultant sliding friction undesirably causes heat build up, and undesirably reduces the mechanical efficiency of the rotary power transmission.\nThe instant inventive continuously variable transmission overcomes or ameliorates the problems discussed above by providing a dual conical friction wheel assembly which is capable of continuously variable power transmission through mechanical deformation of one of the friction wheels of the assembly, such friction wheel comprising an axially skewable scroll.\n(Descriptions below use the spacial orienting terms xe2x80x9cforwardxe2x80x9d, xe2x80x9crearwardxe2x80x9d, xe2x80x9cupperxe2x80x9d, and xe2x80x9clowerxe2x80x9d solely for convenience of reference. The scope of the invention includes all spacial orientations.)\nA primary structural component of the instant inventive continuously variable transmission comprises a rotatable scroll, preferably comprising an axially extending bar axle having a spindle mounted thereover in the manner of a quill assembly, the spindle being capable of rotating about the bar axle, the spindle preferably being further capable of sliding linearly along the bar axle.\nA regular trapezoidally shaped ribbon or strap of sheet steel is preferably wrapped a multiplicity of times around the annular outer surface of the spindle, causing the angled edges of the trapezoidal sheet to form upper and lower conical friction wheel faces.\nWhile the rotatable scroll is symmetrically configured, the pitches of its conical upper and lower faces are equal. Upon an upward skewing of the spiral wrappings of the rotatable scroll to enhance or increase the pitch of its conical lower face, the conical face of its upper face is correspondingly decreased. Conversely, downward skewing of the rotatable scroll increases the conical pitch of its upper face and correspondingly decreases the conical pitch of its lower face. Suitably, the skewable spirally wrapped component of the rotatable scroll may alternately comprise a multiplicity of nesting tubes having progressively narrower axial lengths, such nesting tubes being skewable as described above.\nSecond and third primary structures of the instant invention comprise upper and lower solid or non-skewable conical friction wheels. The conical friction wheels are necessarily held in frictional contact with the upper and lower conical faces of the rotatable scroll. The conical friction wheels are preferably mounted pivotally with respect to the rotatable scroll so that pivoting and counter pivoting motions of the conical friction wheels alternately skew the skewable (preferably spirally wrapped) component of the rotatable scroll upwardly and downwardly, altering the pitches of the scroll\"\"s faces as described above. Upon clockwise pivoting of the upper and lower conical friction wheels with respect to the rotatable scroll, the faces of the scroll are skewed upwardly causing, upon counter rotation of the friction wheels against the scroll, the upper conical friction wheel to rotate faster than the lower conical friction wheel. Conversely, counter clockwise pivoting of the upper and lower conical friction wheels downwardly skews the faces of the scroll, causing the upper conical friction wheel to rotate more slowly than the lower conical friction wheel.\nMounting means in the nature of a support frame is preferably provided, the support frame enabling the upper and lower conical friction wheels to pivot in relation to the rotatable scroll for alternately upwardly and downwardly skewing the faces of the rotatable scroll, the support frame preferably further enabling the upper and lower conical friction wheels to counter rotate against the upper and lower conical faces of the scroll for continuously variable transmission of rotary power through the scroll.\nAccordingly, an object of the instant invention is the transmission of rotary power through a skewable or deformable scroll for continuously variable power transmission.\nOther and further objects, benefits, and advantages of the present invention will become known to those skilled in the art upon review of the Detailed Description which follows, and upon review of the appended drawings."}
{"text": "Adequate effective reservoir characterization is important for effective development of oil and gas reservoirs. Further, the knowledge of formation properties such as permeability, porosity, water saturation etc. may be important for effective well completion design. This may be particularly true for Inflow Control Device (ICD) completion design, intelligent completion design, and horizontal and multilateral well designs. Also, effective reservoir management in an attempt to optimize the production and oil recovery often depends upon good dynamic reservoir characterization. Among all formation properties, permeability (or fluid conductivity) is often the most difficult to determine or predict, and has been a long-standing challenge specifically in the heterogeneous reservoirs having faults, fractures, Vugs or high conductive flow channels. The general perception is the rock with greater porosity usually correlates with greater permeability, which does not hold true in many cases. For example, a formation with small pore throat size may have highly interconnected pores, but these pores and pore channels are too small and the paths available are too restrictive for fluid movement, which substantially reduces the permeability.\nFormation permeability can be directly determined using core plugs in the laboratory, by using a formation fluid tester in a wellbore, or determined with well tests by use of pressure-transient analysis. However, the measured permeability directly from these methods provides permeability at discrete points and has limitations in estimating the permeability continuously across the none-cored wells, particularly in the horizontal wellbore. Also, the amount of core available for direct permeability measurement is typically limited (due to the cost and logistics associated with core sampling), and permeability estimates usually are made by correlations, by use of wireline-log responses. Further, the permeability measured from these sources might have some uncertainty in estimating the permeability at actual reservoir conditions.\nMany investigators have attempted to grasp the complexity of permeability function into a model with general applicability. Most of the models developed are empirical models based on the correlation between formation permeability, formation porosity and irreducible water saturation. These empirical models typically involve measuring porosity and irreducible water saturation of the core and developing mathematical models relating porosity and irreducible water saturation to permeability. In order to use this approach, it is desirable to obtain effective porosity, which is the portion of the porosity that is not isolated and is connected to the pore network and therefore can contribute to fluid flow, and irreducible water saturation. These parameters are not available directly from well logs. Instead, they are estimated from other well log data. However, porosity derived from well log data is not necessarily effective porosity, and methods for deriving irreducible water saturation often rely on effective porosity. Furthermore, empirical models developed for certain formations often perform poorly when used in other fields or formations."}
{"text": "1. Field of the Invention\nThe present invention relates to a thermal head and a thermal printer, and more particularly to a thermal head and a thermal printer using a single crystalline silicon substrate.\n2. Description of the Related Art\nIn recent years, thermal heads for performing thermosensitive recording by selective heat generation of a heat generating element have been used.\nJapanese Patent Application Laid-Open Nos. H02-137943 and H06-320769 disclose a thermal head using a single crystalline silicon substrate.\nThe thermal head disclosed in Japanese Patent Application Laid-Open Nos. H02-137943 and H06-320769 includes: a heat generating element formed on a single crystalline silicon substrate via an insulating film; a driver circuit unit formed on the single crystalline silicon substrate; a wiring layer for connecting the driver circuit unit to the heat generating element; and a protecting film for protecting a thermal head surface.\nA heat generating resistor protrudes toward a printing medium for increasing contact pressure with a printing medium such as thermal paper or an ink sheet to increase printing efficiency.\nProposed methods therefor include a method of partially etching a silicon substrate to form a protruding portion and forming a heat generating portion on the protruding portion, and a method of forming a silicon oxide film having a thickness of several μm on a silicon substrate, then forming a protruding shape on the silicon oxide film by patterning, and forming a heat generating portion on the protruding portion.\nIn the structure in which a heat generating resistor is disposed on the protruding portion formed by partially etching the silicon substrate of the conventional techniques described above, the heat generating resistor is located near the silicon substrate having a large thermal conductivity of 152 W/m·K. Thus, thermal energy generated by the heat generating resistor easily escapes toward the silicon substrate to increase heat loss and thus increase power consumption.\nIn the structure in which the silicon oxide film is formed into a protruding shape, the limit of a height of the protruding portion is several μm due to the limit of stress, and thus sufficient contact pressure with the printing medium cannot be obtained and clear image quality cannot be obtained in some cases.\nThe present invention has an object to provide a thermal head that increases contact pressure between a heat generating portion and a printing medium to increase printing quality with a low heat loss."}
{"text": "Broadband connectivity has revolutionized how individuals and small businesses (collectively referred to as subscribers) connect to the Internet, providing high bandwidth connections with relatively low costs. Broadband connectivity is provided to subscribers in many different forms, including but not limited to: coaxial cable (cable), digital subscriber line (DSL), geosynchronous satellite, wireless, and direct point-to-point optical and wireless. Of course, each method of broadband delivery is usually different and normally is not compatible. Regardless of how the connectivity between the subscriber and the Internet is achieved, there is an underlying commonality that exists in all types of broadband. It is the processing that occurs with data/information packets coming into and out of the subscriber's computer(s).\nThe packets can be divided into one of two groups based on their position in relation to the interface between the subscriber's equipment and the service provider's equipment. Packets can either be local or external. Local packets are packets on the subscriber's side of the interface while external packets are packets on the service provider's side of the interface. Local packets may be packets sent from one computer in the subscriber's network to another in the same subscriber's network (should the subscriber have a computer network) or packets from the subscriber's single computer to itself. When packets are transmitted, processing must be performed on the packets to determine if they are to remain in their respective domains or they are to cross the interface.\nA common solution for processing the local and external packets is via the use of a single processing module that performs needed processing for both local and external packets. The single processing module typically resides at the interface itself. The use of a single processing module permits the sharing of resources and eases design by reducing the overall number of interacting modules. However, the sharing of the single processing module places a large burden on processing resources to perform tasks such as address filtering and translation for local addresses. The address filtering and translation processing may interfere with and affect the performance of the processing of external packets, therefore reducing the overall performance.\nA fairly common solution to the processing of external and local packets is to create two separate modules that are separately responsible for processing external packets and local packets. By separating the processing, it is much less likely that there would be interference since there is no interaction unless a packet crosses the interface. However, the separation of the processes can result in unnecessary duplication of modules, such as the communications protocol modules.\nA need has arisen for method and apparatus that reduces interaction and interference between local and external packet processing and at the same time minimizes the unnecessary duplication of modules."}
{"text": "An HVAC system generally includes a closed loop refrigeration system with at least one evaporator, at least one condenser and at least one compressor. As the refrigerant travels through the evaporator, it absorbs heat from a heat transfer fluid to be cooled and changes from a liquid to a vapor phase. After exiting the evaporator, the refrigerant proceeds to a compressor, then a condenser, then an expansion valve, and back to the evaporator, repeating the refrigeration cycle. The fluid to be cooled (e.g. air) passes through the evaporator in a separate fluid channel and is cooled by the evaporation of the refrigerant. The cooled fluid can then be sent to a distribution system for cooling the spaces to be conditioned, or it can be used for other refrigeration purposes.\nOne type of air conditioner system is a split system where there is an indoor unit or heat exchanger, which is generally the evaporator, and an outdoor unit or heat exchanger, which is generally the condenser. Often, the outdoor unit is placed outdoors and is subject to outdoor ambient conditions, particularly temperature. When the outdoor ambient temperature falls, the amount of heat being removed from the refrigerant in the condenser increases. The increased heat removal in the condenser can result in a decrease in the refrigerant pressure at the suction line to the compressor, commonly referred to as head pressure. The decrease in head pressure results in a lowering of the temperature of the refrigerant at the evaporator. When the temperature of the refrigerant at the evaporator becomes too low, icing of the system can occur. Icing is a condition when the temperature at the exterior of the evaporator is sufficiently low to freeze water present in the atmosphere. The ice formed by the water frozen on the surface reduces the available heat transfer surface and eventually prevents the proper operation of the HVAC system by inhibiting heat transfer and/or damaging system components.\nSome attempts to address the problem of icing have utilized the control of system pressure. In one approach, a variable speed condenser fan or a plurality of condenser fans having independent controls are used to control airflow over the condenser coil. As the amount of air passing over the coil decreases, the amount of heat transfer taking place at the coil decreases. Therefore, the temperature of the refrigerant in the condenser and the pressure of the system increase to allow the indoor coil to cool the air without icing problems. The use of the variable speed condenser fan or a plurality of condenser fans having independent controls has the drawback that it is expensive and requires complicated wiring and controls.\nAn alternate approach for the problem of low system pressure or icing is a parallel set of condensers in the refrigerant cycle, as described in U.S. Pat. No. 3,631,686. The parallel set of refrigerant condensers allows for two modes of operation. One mode of operation allows refrigerant to flow from only one of the refrigerant condensers. During this mode of operation, the condenser that does not permit the flow of refrigerant fills with liquid refrigerant. Because of this flooding, there is a reduction in the effective surface area of the condenser. The reduced surface area thereby reduces the ability of the condenser to remove heat from the refrigerant. Therefore, the temperature of the refrigerant in the condenser and the head pressure of the system increase allowing the indoor coil to cool the air without icing. The use of parallel refrigerant condensers has the drawback that it requires an additional condenser coil and additional piping, thereby increasing the space and cost required for installation. Another drawback associated with refrigerant flooding of the condenser coil is the resultant decrease in system capacity. Refrigerant normally available in a properly operating system is trapped in the condenser coil and not available to the compressor, thereby decreasing system capacity.\nAn additional alternate approach for the problem of low system pressure is the use of a valve that controls the discharge or flow of liquid refrigerant from the condenser to a receiver vessel downstream of the condenser to maintain control of the amount of condensing surface exposed to the outside temperature, as described in U.S. Pat. No. 2,874,550. The discharge of refrigerant from the condenser is controlled by a pressure-response valve that mechanically opens to allow the flow of liquid refrigerant from the condenser to the receiver vessel reducing the level of liquid inside the condenser, thereby lowering the system pressure. Alternatively, the valve is closed to stop the flow until the level of refrigerant rises in the condenser in an amount that reduces the effective cooling surface of the condenser. The reduced surface area thereby reduces the ability of the condenser to remove heat from the refrigerant, thereby raising the pressure of the system. The use of a pressure-response valve and a vessel downstream of the condenser to maintain control of the amount of condensing surface has the drawback that it includes a specially designed valve and additional piping, thereby increasing the required space and cost. As discussed above, another one of the drawbacks with refrigerant flooding the condenser coil is decreased system capacity. Refrigerant normally available in a properly operating system is trapped in the condenser coil and not available to the compressor, thereby decreasing system capacity.\nAn additional alternate approach for the problem of low system pressure is the use of a refrigerant bypass around the condenser, as described in U.S. Pat. No. 3,060,699 and U.S. Reissued Pat. No. Re. 27,522. If the temperature and pressure of the refrigerant in the condenser are sufficiently high, a valve will close on a condenser bypass and the flow of refrigerant will be directed to the condenser. If the temperature and pressure of the condenser are not sufficiently high, the valve will open on a condenser bypass and at least some of the flow of refrigerant will be directed away from the condenser. The result of the bypass is an increase in pressure through the pipe leading to the evaporator downstream of the compressor. The use of a bypass has the drawback that it includes a specially designed valve and additional piping, thereby increasing the required space and cost.\nWhat is needed is a method and system for controlling the system refrigerant pressure without the drawbacks discussed above."}
{"text": "The invention relates to a fuel supply system for a spark-ignition internal combustion engine having arranged on each cylinder a blow-in nozzle for the introduction of a mixture of fuel and a combustion gas into a combustion chamber.\nSpark-ignition internal combustion engines have an unsatisfactory efficiency in particular in the part-load range as a result of high charge exchange losses. One possibility for reducing said charge exchange losses and thereby improving efficiency resides in exhaust gas recirculation. By means of exhaust gas recirculation, the engine can operate unthrottled during part-load operation, which reduces the required charge exchange work of the engine and increases the compression pressure.\nAn increase in the proportion of recirculated exhaust gas can however also bring about a degradation of the ignitability and combustion of the mixture of fuel, air and exhaust gas. Misfiring can even occur.\nIn order to utilize the advantages of a high exhaust gas recirculation rate and in order to obtain the ignitability required for this purpose, DE 199 45 544 A1 discloses a fuel supply system for a spark-ignition internal combustion engine in which a first fuel metering device comprises an intake pipe injector which is arranged in an intake duct which opens out into the combustion chamber. By means of the intake pipe injector, a fuel quantity is injected in order to form a lean mixture with a predefined exhaust gas proportion. For reliable ignition, one blow-in valve is provided per cylinder, which blow-in valve, in order to ignite the lean mixture, blows a mixture of a small liquid fuel quantity and air as a combustion gas into the combustion chamber. The quantity which is blown in can be easily and reliably ignited. After ignition, the mixture quantity ignites the remaining mixture.\nIt is the principal object of the present invention to provide an internal combustion engine with a fuel supply system which is small and effectively supplies a mixture of fuel and combustion air into the cylinders of the engine."}
{"text": "This application is based upon and claims benefit of priority of Japanese Patent Application No. Hei-11-366454 filed on Dec. 24, 1999, the content of which is incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to a method of manufacturing a semiconductor device on an SOI (Silicon on Insulator) substrate, and more particularly to a method of forming an oxidized film on a silicon layer of the SOI substrate.\n2. Description of Related Art\nIn a conventional method of manufacturing capacitors or MOS transistors on a silicon substrate, an impurity is diffused into the silicon substrate in a high density to form a highly doped layer, and then an oxidized film is formed on the highly doped layer. Electrodes made of poly-silicon or the like are formed on the silicon dioxide film.\nIt is frequently required that such a silicon dioxide film has a high quality to be durable against a high voltage. However, such a high quality is difficult to be realized when the oxidized film is formed on the highly doped layer. This is because the highly doped layer tends to become a gettering site and the silicon dioxide film is contaminated with contaminants such as heavy metals, and the durability of the silicon dioxide film is adversely affected by the contaminants.\nIt is possible to form an additional gettering layer in the silicon substrate to improve the quality of the oxidized film. However, it is generally difficult to form such a gettering layer in the SOI substrate. JP-A-10-25621 proposes to form the gettering layer in the SOI substrate, but at least an additional costly process to form the same is required in a series of manufacturing processes of a semiconductor device.\nThe present invention has been made in view of the above-mentioned problem, and an object of the present invention is to provide a method of forming a quality oxidized film on the SOI substrate without using the process of forming the gettering layer.\nAn impurity such as boron or phosphorus is diffused into a silicon layer on an SCI substrate by ion-implantation and driven-in, thereby forming a diffused layer having a high density of the impurity on the SOI substrate. Then, an oxidized film is formed on the diffused layer by thermal oxidation. The oxidized film is heat-treated at a high temperature of 1100xc2x0 C. or higher, preferably 1150xc2x0 C. or higher. When boron is diffused to form the diffused layer, the high temperature heat treatment is performed for 20 minutes or longer. When phosphorus is diffused, the high temperature heat treatment is performed for a period of 20 minutes or longer but not exceeding 60 minutes.\nBy applying the high temperature heat treatment to the oxidized film, contamination in the oxidized film is driven-out, or annealed-out. Thus, the quality of the oxidized film such as durability against a high voltage is improved.\nTo form a capacitor using the oxidized film, electrodes made of a material such as poly-silicon are patterned on the oxidized film. Other semiconductor elements such as an LDMOS (Lateral Diffused MOS) may be formed on the same SOI substrate in addition to the capacitor. In this case, the high temperature heat treatment can be performed as a common process for improving the quality of the oxidized film and for forming other semiconductor elements.\nFurther, plural capacitors, e.g., a first capacitor and a second capacitor may be formed using the oxidized films having respectively different thickness. In this case, a preliminary oxidized film covering both areas for the first and second capacitors is first formed. Then, only the oxidized film covering the second capacitor area is removed by wet-etching, keeping the preliminary oxidized film for the first capacitor. Then, the silicon layer is again thermally oxidized to form a new oxidized film covering both areas for the first and second capacitors. Thus, a thin film is formed for the second capacitor and a thick film is formed for the first capacitor. The high temperature heat treatment is performed either immediately after both oxidized films are formed or after the capacitor electrodes are formed on the oxidized films. The quality of both oxidized films is improved by a single high temperature heat treatment.\nAccording to the present invention, the quality of the oxidized film or films formed on the SOI substrate is improved by the high temperature heat treatment without forming an additional gettering layer.\nOther objects and features of the present invention will become more readily apparent from a better understanding of the preferred embodiment described below with reference to the following drawings."}
{"text": "A lithographic apparatus is a machine that applies a desired pattern onto a target portion of a substrate. Lithographic apparatus can be used, for example, in the manufacture of integrated circuits (ICs). In that circumstance, a mask, which is alternatively referred to as a mask or a reticle, may be used to generate a circuit pattern corresponding to an individual layer of the IC, and this pattern can be imaged onto a target portion (e.g. comprising part of, one or several dies) on a substrate (e.g. a silicon wafer) that has a layer of radiation-sensitive material (resist). In general, a single substrate will contain a network of adjacent target portions that are successively exposed. Known lithographic apparatus include so-called steppers, in which each target portion is irradiated by exposing an entire pattern onto the target portion in one go, and so-called scanners, in which each target portion is irradiated by scanning the pattern through the beam in a given direction (the “scanning”-direction) while synchronously scanning the substrate parallel or anti parallel to this direction.\nProgressive improvements of a lithographic apparatus have been made which have allowed patterns with smaller resolutions to be projected onto substrates. One such improvement involves providing a liquid between a projection system of a lithographic apparatus and a substrate. This provides the projection system with a numerical aperture (NA) greater than 1.0 (e.g. 1.35NA)."}
{"text": "Electronic devices, such as portable electronic devices, have gained worldwide popularity due to their broad applications. Portable electronic devices may include, for example, smartphones, wireless personal digital assistants (PDAs), tablets, laptop computers with wireless or Bluetooth® capabilities, cellular telephones, etc. A user may use, for example, a smartphone to perform a variety of functions including making telephone calls, sending electronic messages, taking photos, reading articles, and other functions, by installing applications.\nWith the increasing number of applications, users tend to install more and more applications on their portable electronic devices. At the same time, users have started reading articles and/or news on their portable electronic devices. To accommodate users' various usages, manufacturers have brought into the market portable electronic devices that have relatively large touchscreens to enable rendering of more icons and contents. This, however, gives rise to problems for users who operate portable electronic devices by using one hand. For example, if a user operates a portable electronic device by one hand, and touches the touchscreen of the portable electronic device with a thumb, the thumb of the user may not reach icons that are located far from the thumb location on the touchscreen. In this case, the user may need to instead use both hands to operate the portable electronic device via the touchscreen.\nTo solve this problem, existing technologies relating to Human-Computer Interaction may provide a settings menu, by which a user sets his/her hand use preference. For example, if a user prefers to use both hands to operate a touchscreen, the user can set a corresponding preference on the settings menu.\nHowever, the existing technologies are not only complicated but also require presetting. Situations may arise in which a user has a need to use both hands, even though the user usually uses one hand to operate a portable electronic device. As an example, when the user is reading an article on the portable electronic device while being a passenger in a vehicle, the user may need to use both hands to avoid shaking or dropping the portable electronic device. Due to this need, the user may not have time to preset a user preference to use both hands. As a result, a user may drop a portable electronic device while only one hand is being used to hold the portable electronic device."}
{"text": "The present invention is related to an apparatus for controlling the indexing of hydraulic stepping motor and in particular to a servo-valve for controlling the flow of htdraulic fluid thereto.\nHydraulic stepping motors are known which are constructed as a linear motor in form of a pressure fluid cylinder or in form of a hydraulic rotating motor. In either case it is known to combine such motors with a control or servo-valve comprising slide valve having a reciprocable linear shift member or with a member reciprocably rotatable about a central axis (see German Pat. No. OS 1,550,365). In connection with slide valves having a linear shiftable member it is possible and known to use also asymmetric pistons.\nIt is an object of the instant invention to provide improved apparatus for the control of hydraulic motors which may use any of these combinations. The apparatus may also be provided known intermediate amplifiers or gears altering the gear ratio between the hydraulic motor and the control valve.\nIt is the object of the instant invention to develop a hydraulic step motor of the type described in such a way as to render possible in a simple construction a high number of step positions and one having little leakage. In addition it is the object to provide a device having a very good dynamic performance, an insensitiveness against varying load moments, higher start-stop-frequencies and step losses.\nAccording to the invention these tasks are solved by providing in conjunction with a hydraulic stepping motor a servo-control system comprising a slide valve comprising a housing and a slidable valve member arranged in, and coextensive with, the housing. The slidable valve member includes at least two fluid manifolds conjointly movable therewith and a plurality of primary positioning webs and secondary sealing webs alternately spaced from one another and uniformly distributed over its periphery. These projections define a plurality of intermediate gaps alternately connected with one or the other of the fluid manifolds. A plurality of uniformly distributed channels provided in the housing opens against these webs. The channels have a width identical with the peripheral width of the primary projections, but have mutual distances different from those existing between the primary projections and are fewer in number than these primary projections.\nThe main connecting conduits provide communication between the source of fluid under pressure and the fluid manifolds. A plurality of return conduits, each of which is connected to one of the channels, is coupled to a selector valve comprising a plurality of control valves each communicating with at least one of the return conduits. The control valves are actuable to couple a selected return conduit to a return duct outlet opening into a sump whereto the source of fluid under pressure is connected to receive therefrom the fluid it pressurizes.\nIn idle condition a uniform pressure exists as long as the valves in the selector are closed. In response to the opening of one of the valves, the corresponding intermediate gap discharges the fluid through its corresponding channel and the return fluid conduit, via the opened valve into the sump, thus a drop of pressure being provoked in the discharging gap and the step motor being caused to revolve in the decreased pressure direction. Then, one of the primary webs moves towards and covers the previously opened channel and finally blocks it. Therefore, the rotary valve member ceases to move and the fluid escape channel closure restores the pressure in the servo-valve, so that the rotation stops after a predetermined step only. Thus an effective control of step rotation is provided.\nPreferably the servo-valve is of a rotary type including a rotatable cylindrical member housed in a cylindrical wall. The webs comprise radial projections. However, linear housing and slide members may be used since these too can be made with projecting webs and channels.\nAn advantage of the foregoing construction lies in the fact that by connecting the gaps between the webs, on the same facial sides, to a common manifold, indexing in opposite directions can be obtained in one or more steps.\nIt is particularly useful if in the preferable embodiment each separating web is broadened up to form a sealing member in such a way that in each position of the movable member at least some of the channels in the valve housing are covered by the sealing web.\nThe present invention results in a high number of indexing positions, together with a very small possibility of leakage. Due to this diminished leakage the dimensions of the slide valve may be significantly decreased in comparison with known systems having an equal number of positions. Furthermore, the number of control valves necessary in the selector valve device may be substantially decreased so that there is a remarkable decrease in constructive expenditure. The width of the separating webs, respectively, may in order to obtain a favorable behavior in operation, be chosen so that at both sides of each positioning web there remain gaps which are slit like ducts having a width corresponding to the width of the positioning web itself.\nPreferably, the number of channels is nearly equal to the number of the positioning webs. Most favorably the number of the positioning webs and the number of the channels differ only to the extent of the valve 1.\nFurthermore, it is useful if the conduits of two or more non-adjacent channels are connected to the return conduit via a single control valve of the selector valve device.\nThe invention will be described in more detail in the following with reference to the attached drawings showing two embodiments."}
{"text": "The present invention relates to a shift interlock device, and more particularly to a brake-transmission-shift-ignition (BTSI) interlock device and circuit integrated into a shifter.\nBrake-transmission-shift-ignition (BTSI) interlock devices and circuits are known in the art, and further it is known to incorporate a toggle mechanism and preassembled switch into such devices. For example, see U.S. Pat. No. 5,759,132 to Osborn, issued Jun. 2, 1998. A known existing BTSI shifter similar to U.S. Pat. No. 5,759,132 is described in the discussion below entitled xe2x80x9cPrior Artxe2x80x9d. (See FIGS. 1-2 of the present disclosure.) This known existing BTSI shifter includes a preassembled micro switch that is soldered onto a circuit board. The subassembly is then operably positioned in and secured to the handle of a shift lever assembly of a vehicle shifter. However, this shift lever assembly is undesirably expensive, partially because of the cost of the circuit board, but also because of secondary process costs (e.g. soldering, manual placement and attachment within the handle, and electrical connections). Also, circuit boards have quality and warranty concerns because they are not well suited for the vibrations and harsh environments commonly associated with vehicle shifters (e.g. temperature and humidity variations in the passenger compartment associated with day and night, and/or temperature extremes associated with winter and summer, and/or vibrations and with shifting the shift lever and/or associated with traveling at high speeds along a bumpy road and/or when an engine idles). Accordingly, further improvement is desired in this assembly, including reducing component costs and increasing the automation and efficiency of assembly, and improving the durability and robustness of the shift lever assembly.\nIt is also known to use a lead frame in a brake shifter interlock having a toggle mechanism. For example, see Withey U.S. Pat. No. 5,938,562, which discloses a lead frame incorporated into a toggle interlock device. In the Withey arrangement, the lead frame incorporates conductive components forming a switch. But this switch arrangement can have quality problems, since the conductive components are subject to distortions and dimensional variations, wear, and other problems that occur during installation and use and during the wide temperature variations commonly experienced by shifters. This can lead to poor and unreliable operation of the integral switch. It is desirable to incorporate a preassembled switch into a circuit using a lead frame. However, problems still remain in terms of assembly and warranty problems associated with soldering and/or other electrical connections. Also, the problems associated with dimensional inconsistencies and part-handling common in lead frames need to be addressed, as well as the overall ability to automatically assemble the components.\nAccordingly, a shifter is desired solving the aforementioned problems and having the aforementioned advantages.\nThe present invention includes a shifter for a vehicle having an electrical control circuit. The shifter includes a base, a shift lever pivoted to the base, and an interlock device on one of the base or the shifter that engages an abutment surface on the other. The abutment surface is configured to be selectively engaged by the interlock device to control movement of the shift lever. The interlock device includes a preassembled switch, an electromechanical device, and a lead frame having at least four conductors. The electromechanical device has an interlock member movable to an extended position for engaging the abutment surface and movable to a retracted position for operating the switch and allowing the shift lever to move. The four conductors operably interconnect the switch and the electromechanical device and define a three-prong terminal adapted for electrical connection to the vehicle control circuit for operating the electromechanical device and for signaling to the vehicle control circuit that the interlock member has been operated.\nIn another aspect of the present invention, a preassembled interlock device includes a housing and a toggle interlock mechanism including an extendable pin operably positioned in the housing. The pin is extendable to a position outside the housing for engaging an abutment surface. A preassembled switch, an electromechanical device, and a lead frame are attached to the housing in an arrangement where the electromechanical device operates the switch when the extendable pin is retracted. The lead frame includes at least four conductors operably interconnected to the switch, the electromechanical device, and to a terminal adapted for electrical connection to a control circuit for operating the electrical mechanical device, and for signaling to the vehicle control circuit that the extendable pin has been moved.\nIn still another aspect of the present invention, a method of assembly for an interlock device comprises steps of providing a housing, and positioning a lead frame in the housing. The method further includes operably positioning a toggle interlock mechanism in the housing, the interlock mechanism including an electromechanical device having a pin extendable to a position outside of the housing. The method also includes positioning a preassembled switch in the housing, including electrically connecting the switch to the lead frame. The method also includes electrically connecting the electromechanical device to the lead frame so that the electromechanical device operates the switch when the extendable pin is retracted. The method further includes separating parts of the lead frame to form at least four separate conductors that operably interconnect the switch and the electromechanical device to a terminal adapted for electrical connection to a control circuit for operating the electromechanical device, and for signaling to the control circuit that the interlock member has been extended.\nIn yet another aspect of the present invention, a preassembled interlock device includes a housing having a plurality of protrusions, and a toggle interlock mechanism including an extendable pin positioned in the housing but extendable to a position outside of the housing. The interlock device also includes an electromechanical device attached to the housing, a preassembled switch in the housing including a plurality of first contacts, and a lead frame having at least four conductors including a plurality of second contacts. When the electromechanical device is energized, it retracts the pin and operates the switch. The first contacts are positioned and oriented to telescopingly engage the second contacts during assembly. The second contacts each include a mechanical retainer that retains the first contacts in operable engagement with the second contacts and further they are adapted to electrically operably interconnect the electrical mechanical device to a control circuit for operating the electrical mechanical device. The lead frame further includes a plurality of location holes with retaining tines that engage the housing protrusions to non-releasably and accurately locate the lead frame in the housing.\nIn still another aspect of the present invention, a shifter for a vehicle having an electrical control circuit includes a base component, a shift lever component operably positioned on the base component, and an interlock device on one of the base and shift lever components. The interlock device is configured and adapted to selectively engage an abutment surface on the other of the components for interlocking the shift lever in a selected gear position. The interlock device includes a preassembled switch and a lead frame, with the preassembled switch having at least three first contacts and the lead frame having at least three second contacts engaged with the first contacts. The first contacts are positioned and oriented to telescopingly engage the three second contacts during assembly, and at least one of the first and second contacts each include a mechanical retainer for securely retaining the first contacts in operable engagement with the second contacts after assembly.\nThese and other aspects, objects, and features of the present invention will be understood and appreciated by those skilled in the art upon studying the following specification, claims, and appended drawings."}
{"text": "The present invention relates to a furnace and a method for drawing an optical fiber from heated and molten preform.\nVarious techniques have been known relating to a furnace for drawing an optical fiber from heated preform formed with silica glass as a main component, and they are described in Japanese Patent No. 2,542,679, Japanese Patent Laid-Open No. 147969/1993 and Japanese Patent Laid-Open No. 2832/1997. Since the figures and terminology in these prior arts differ partially from those used in this invention, descriptions will be made below where they are interpreted according to the figures and terminology used in this invention in order to clarify the difference between the prior arts and this invention.\nA main part of the furnace for drawing an optical fiber disclosed in Japanese Patent No. 2,542,679 is shown in FIG. 8, in which numeral 21 denotes a preform, 21a denotes an optical fiber, 22 denotes an inner tube, 22a denotes a gas blowing inlet, 23 denotes an outer tube, 24 denotes a gas supplying inlet, 24a denotes a gas passage, 25 denotes a dummy rod, 25a denotes a connecting part, 26 denotes a retainer, 27 denotes a seal piston, 28 denotes a muffle tube, and 29 denotes a heater.\nIn the furnace for drawing an optical fiber, the dummy rod 25 and the preform 21 are arranged inside the muffle tube 28 and the inner tube 22 arranged to be connected to an upper end of the muffle tube 28, while they are connected via the connecting part 25a and are descended together. The vicinity of a lower end of the preform 21 is melted through heating by the heater 29 arranged outside the muffle tube 28, and the optical fiber 21a is drawn downward from the lower end of the preform 21. The inner tube 22 arranged to be connected to the upper end of the muffle tube 28 is to contain the long preform 21 on starting the drawing.\nThe gas passage 24a is formed between the inner tube 22 and the outer tube 23 arranged outside the same, and an inert gas is supplied from the gas supplying inlet 24 to the gas passage 24a, so as to blow the inert gas into the inside of the inner tube 22 from the numerous gas blowing inlets 22a provided circumferentially and in the height direction on the wall surface of the inner tube 22. The inert gas is flowed inside the inner tube 22 and the muffle tube 28 to prevent oxidation deterioration of the muffle tube and the like, and when the temperature distribution of the inert gas by heating and the flow of the inert as are not uniform, the fluctuation of the diameter of the optical fiber drawn from the preform is liable to occur.\nTherefore, in this example of the furnace for drawing an optical fiber, the seal piston 27 connected to the dummy rod 25 via the retainer 26 and moved with the dummy rod 25 is provided at the dummy rod 25 arranged on the upper part of the preform 21. In the beginning of drawing, the dummy rod 25 and the seal piston 27 are in the upper part because the preform is long. With proceeding the drawing, the preform 21 is shortened from the lower end, and is descended, and thus, the dummy rod 25 and the seal piston 27 are also descended.\nIn this case, if the seal piston 27 were not present, the space between the dummy rod 25 and the inner tube 22 would gradually increase, but because the seal piston 27 is present, the volume of space above the preform 21 is substantially constant. Therefore, it has been said that turbulence of the stream of the inert gas in the space between the preform 21 and the seal piston 27 occurs only scarcely due to the provision of the seal piston 27.\nIn the furnace using the seal piston for drawing an optical fiber, when the preform has a length of 1.5 m or more, the seal piston necessarily has an proportionate length, and the weight thereof becomes also heavy. Because a supporting member for supporting them at the upper part must withstand the weight of the preform and the seal piston, the supporting member becomes also necessarily large. Because the seal piston must resist to high temperature, a heat resistant material, such as carbon, quartz and the like, is necessarily used for the seal piston, and it becomes costly when it is of large scale.\nFurthermore, because the seal piston moves as sliding on the inner wall surface of the inner tube, dust is liable to come from the sliding part, which may adversely affect the strength of the drawn optical fiber.\nFurthermore, as the seal piston descents, the numerous gas blowing inlets provided in the wall surface of the inner tube are sealed one by one from the upper part with the seal piston, a precise controller for continuously controlling the gas flow rate is necessary to maintain the constant flow rate of the stream of the inert gas.\nA furnace for drawing an optical fiber disclosed in Japanese Patent Laid-Open No. 147969/1993 will then be described. A main part of the furnace for drawing an optical fiber is shown in FIG. 9. In FIG. 9, the same references as in FIG. 8 show the same components. Numeral 30 denotes a separating plate, 30a denotes a gap, 30b denotes pores, 31a denotes an upper space, 31b denotes a lower space, and 32 denotes an upper lid. The furnace for drawing an optical fiber shown in FIG. 9 is different from that shown in FIG. 8 in the following points. There is no member corresponding to the seal piston in FIG. 8, and an upper end of an inner tube 22 is closed with the upper lid 32 except for the part through which a dummy rod 25 penetrates.\nIn the furnace for drawing an optical fiber, a preform 21 is arranged inside a muffle tube 28 and the inner tube 22 connected to an upper end thereof, and the preform 21 is supported by hanging by the dummy rod 25 through a connecting part 25a. The vicinity of the lower end of the preform 21 is heated and melted by a heater 29 arranged outside the muffle tube 28, and an optical fiber 21a is drawn downward.\nAn outer tube 23 is arranged concentrically outside the inner tube 22, and an inert gas is blown into the inside of the inner tube via a gas blowing inlet 22a through a gas passage 24a formed with the outer tube 23 and the inner tube 22. The inert gas is introduced to the gas passage 24a from a gas supplying inlet 24. The inert gas blown into the inner tube 22 descends through the space between (the inner tube 22 and the muffle tube 28) and (the preform 21 or the dummy rod 25), then it is exhausted through the vicinity of the optical fiber 21a. \nThe space inside the inner tube 22 is separated by the separating plate 30 comprising a quartz plate or the like into an upper part and a lower part, and the inert gas flows from the upper space 31a of the separating plate 30 to the lower space 31b through the gap 30a between the separating plate 30 and the inner tube 22 or the pores 30b provided in the separating plate 30. Because the lower space 31b is of a relatively high temperature and turbulence of the gas stream is decreased owing to the presence of the separating plate 30, it has been said that an optical fiber having a small fluctuation in diameter can be drawn even in the case of a large preform.\nHowever, in the case of this furnace for drawing an optical fiber, when the preform 21 becomes small with proceeding the drawing, the upper space 31a becomes large. On the other hand, the temperature near the upper end of the preform 21 is increased by heating the vicinity of the lower end of the preform 21, and said temperature becomes higher when the preform becomes small. Because of the presence of the separating plate 30 comprising a quartz plate, the temperature is somewhat decreased above the separating plate, but it becomes 550xc2x0 C. or more near the lower end of the upper space 31a. At this time, the temperature near the upper end of the upper space 31a is about 200xc2x0 C. to form a considerable difference in temperature between the upper part and the lower part inside the upper space 31a, therefore convection is liable to be caused. The convection causes fluctuation in volume of the inert gas, and the fluctuation is transmitted to the lower space 31b by the flow of the inert gas through the gap 30a or the pore 30b of the separating plate 30. Thus, fluctuation in flow of the inert gas near the optical fiber occurs. As a result, the amount of heat transmission is changed by the fluctuation of the inert gas to easily cause fluctuation in viscosity and softened amount of glass, and thus it is difficult to suppress fluctuation in fiber diameter to a small value less than the prescribed value.\nA furnace for drawing an optical fiber disclosed in Japanese Patent Laid-Open No. 2832/1997 will then be described. The furnace for drawing an optical fiber is shown in FIG. 10. In FIG. 10, the corresponding references as in FIG. 8 show the same components. Numeral 33 denotes an upper space above the preform, and 34 denotes an auxiliary heater. In the furnace for drawing an optical fiber, a preform 21 is arranged inside a muffle tube 28 and the inner tube 22 connected to an upper end thereof, being supported by hanging by the dummy rod 25, and the vicinity of a lower end of the preform 21 is heated and melted by a heater 29 arranged outside the muffle tube 28, so as to draw an optical fiber 21a downward.\nAn outer tube 23 is arranged outside the inner tube 22 concentrically, and an inert gas is blown into the inside of the inner tube via a gas blowing inlet 22a through a gas passage 24a formed with the outer tube 23 and the inner tube 22. Because the vicinity of the upper end of the preform 21 becomes a high temperature of 1,000xc2x0 C. or more by heating the vicinity of the lower end of the preform 21, a large temperature difference formed between the upper part and the lower part of the space 33 above the preform.\nIn order to decrease the temperature difference, the auxiliary heater 34 is arranged outside the upper end of the inner tube 22 to heat the vicinity of the upper end of the inner tube 22. Thus, the vertical temperature difference of the space 33 between the inner tube 22 and the dummy rod 25 above the preform is decreased by the heating by the auxiliary heater 34 to prevent the convection from occurring at the space 33, and fluctuation in diameter of the optical fiber 21a drawn from the preform 21 is decreased.\nHowever, in the case of this furnace for drawing an optical fiber, because the space 33 above the preform becomes large as the drawing of a large scale preform proceeds, it is difficult to make the temperature uniform inside the space solely by an auxiliary heater 34 arranged outside the upper part of the inner tube 22.\nIt is then considered that several auxiliary heaters are arranged vertically not only at the upper part but also at the lower part of the inner tube, and by controlling the temperature of each auxiliary heater, the temperature of the space above the preform is made uniform. However, there are problems that cost of the apparatus becomes increased, and the temperature control becomes complicated.\nThe invention is to provide a furnace for drawing an optical fiber and a method for drawing an optical fiber to solve the problems associated with the conventional art described in the foregoing.\nThe furnace for drawing an optical fiber according to the invention comprises a muffle tube and an inner tube connected to an end of the muffle tube. The preform is arranged inside the muffle tube and the inner tube and supported by a dummy rod at an upper part thereof in such a manner that the preform descends with the dummy rod. The preform is heated from outside of the muffle tube by a heater, and melted so as to draw an optical fiber from a lower end of the preform. One set or plural sets of separating plates separating a space inside the inner tube above the preform into plural parts in an advancing direction of the preform are arranged inside the space, and a gas blowing inlet for blowing an inert gas into the inner tube and the muffle tube is provided at the wall of the inner tube at a part under the separating plates.\nTherefore, the inert gas entered from the gas blowing inlet mainly flows downward, and thus the gas in the upper space of the separating plate above the preform less influences the conditions of the drawing of the preform. Accordingly, the drawing is stably continued, and fluctuation in fiber diameter is decreased. Therefore, an optical fiber produced by using this furnace for drawing an optical fiber has a substantially constant outer diameter in the longitudinal direction, and has small fluctuation in transmission characteristics. Owing to the use of a simple separating plate but not using a large member such as a seal piston, the cost of the apparatus is low, and since heavy sliding as in the seal piston does not occur, dust formed by sliding is so minimal that it doesn\"\"t adversely affect the strength of the optical fiber.\nWhile the separating plates comprise plural sets of separating plates arranged as being penetrated by the dummy rod, and the plural sets of separating plates descend with the descent of the dummy rod, the respective sets of the plural sets of separating plates are stopped on the inner wall surface of the inner tube one by one, so that the space inside the inner tube above the preform is respectively separated into an upper part and a lower part by the stopped separating plates. Thus, because the space inside the inner tube can be separated by plural positions into suitable sizes, the flow of the inert gas inside the inner tube can be further stabilized even when the preform becomes small and the space inside the inner tube is enlarged with proceeding the drawing of the optical fiber.\nWhile the respective sets of separating plates may each comprise one disk-shaped plate, the sets of separating plates may each be constituted by two plate members of an outer member and an inner member. In this case, the outer member has an outer diameter that is the same as the inner diameter of the inner tube at a position where the separating plate is stopped, and has a center hole diameter that is larger than the outer diameter of dummy rod so as to absorb the deviation from the concentric condition of the inner tube and the dummy rod. The inner member has an outer diameter that is larger than the center hole diameter of the outer member and smaller than the outer diameter of the outer member, and has a center hole diameter larger than an outer diameter of the dummy rod by 2 to 10 mm.\nWhen the dummy rod penetrates both the members with the inner member being placed at upper side and the outer member being placed at lower side, the upward inner member is supported by being placed on the outer member. By constituting each of the sets of separating plates of an outer member and an inner member, even when the concentric condition of the inner tube and the dummy rod is deviated by rolling of the preform to cause decentering, the upper member moves in the radial direction by sliding on the lower member, and no such phenomenon that the inner wall surface of the inner tube is damaged by the separating plate occurs.\nAlternatively, one or plural sets of separating plates may be arranged near the lower end of the dummy rod or at the upper part of the preform by fixing to the dummy rod, the connecting part or the preform, and they descend with the descent of the preform without stopping on the way on the inner wall of the inner tube. Accordingly, the space between the separating plates and the preform can be unchanged and constant even when the preform is shortened, and thus the flow of the gas in the space above the preform can be stabilized. It is also possible in this case that the respective sets of separating plates are each constituted by an outer member and an inner member. In this case, however, the outer diameter of the outer member must be smaller than the inner diameter of the inner tube by 5 to 10 mm, and the center hole diameter of the outer member must be larger than the outer diameter of the dummy rod. The outer diameter of the inner member must be larger than the center hole diameter of the outer member but smaller than the outer diameter of the outer member. The center hole diameter of the inner member must be the same as or larger than the outer diameter of the dummy rod. The inner member is then supported by the dummy rod or the connecting part with the inner member being downward and the outer member being upward.\nBy constituting each of the respective sets of separating plates by the outer member and the inner member, even in the case where the concentric condition is deviated by rolling of the preform to cause decentering, the upper member moves in the radial direction by sliding on the lower member, and no such phenomenon that the inner wall surface of the inner tube is damaged by the separating plate occurs.\nWhen plural protrusions are provided on an outer periphery of the separating plate, the sliding friction coefficient of the separating plate to the inner wall surface of the inner tube on descending can be reduced, so as not to damage the inner tube, because only the protrusions are in contact with the inner wall surface of the inner tube.\nBy forming the separating plate with a heat insulating material, heat transferred from the upper end of the heated preform to the upper space above the separating plate can be reduced, and thus the temperature of the upper space above the separating plate can be reduced. Therefore, the temperature difference between the upper part and the lower part in the upper space is also reduced, and the generation of convection of the inert gas due to the temperature difference can also be suppressed. As a result, the influence on fluctuation in gas flow near the lower part of the preform influenced by the convection of the inert gas in the upper space can be suppressed, and thereby fluctuation in heat transfer amount by the inert gas, and fluctuation in viscosity and softened amount of the glass are suppressed, so as to reduce fluctuation in diameter of the optical fiber. Furthermore, because the influence by the convection of the inert gas in the upper space is reduced, the inert gas atmosphere can be maintained and the fluctuation in fiber diameter can be suppressed by flowing only a small amount of the inert gas near the preform. Therefore, the amount of consumption of high cost inert gas, such as helium, is reduced, and thus an economical effect is also expected.\nFurthermore, by arranging an auxiliary heater outside the vicinity of the upper end of the inner tube to heat the vicinity of the upper part of the inner tube, the temperature inside the upper space can be further uniform."}
{"text": "It is well known that a mixture of fluids having different densities may be separated from one another through use of a centrifugal separator. One specific use of such a separator is in the separation of oil from gas vented from a crank case forming part of an internal combustion engine.\nWith regard to this specific use of separators, there can be a tendency for the high pressure gasses found in the combustion chambers of an internal combustion engine to leak past the associated piston rings and into the crank casing of the engine. This continuous leaking of gas into the crank case can lead to an undesirable increase of pressure within the crank case and, as a consequence, to a need to vent gas from said casing. In large commercial vehicles, vented gas is generally reintroduced into the inlet manifold of the engine. However, the gas vented from the crank casing typically carries a quantity of engine oil (as droplets or a fine mist), which is picked up from the reservoir of oil held in the crank casing. More specifically, gas flowing between an engine cylinder and the associated piston tends to pick up lubricating oil located on the cylinder wall. Also, condensing of oil vapour by an engine's cylinder block cooling system generates an oil mist in the crank case.\nIn order to allow vented gas to be introduced into the inlet system without also introducing unwanted oil (particularly into a turbocharging system wherein the efficiency of the compressor can be adversely affected by the presence of oil), it is necessary to clean the vented gas (i.e. to remove the oil carried by the gas) prior to the gas being introduced into the inlet system. This cleaning process may be undertaken by a centrifugal separator, which is mounted on or adjacent the crank case and which directs cleaned gas to the inlet system and directs separated oil back to the crank case.\nThere are a number of problems associated with some prior art ALFDEX™ separators. These problems can be considered in three broad categories.\nFirst, the fluid pathways through the separator give rise to pressure losses which adversely affect the flow capacity of the separator and, consequently, the size of engine with which the separator can be used.\nSecond, the arrangement of some of these prior art separators is such that, under certain conditions, cleaned gas can become contaminated before leaving the separator.\nThird, certain manufacturing techniques and construction features associated with these prior art separators can lead to assembly difficulties and/or reliability problems."}
{"text": "This invention relates to a process for the melt-spinning of a filamentary structure from a synthetic polyamide polymer. More particularly, it is concerned with a continuous process for the formation of high strength polycaproamide filament, yarn or the like by melt-spinning a synthetic linear fiber-forming polycaproamide having excess number of carboxyl end groups over amino end groups.\nU.S. Pat. No. 2,174,527 discloses that dibasic dicarboxylic acids when employed in excess in a diacid-diamine polyamide, serve to terminate the polymer, thereby minimizing further increase in polymer molecular weight.\nU.S. Pat. No. 2,241,321 teaches polymerization of caprolactam in the presence of a diamine, followed by addition of sebacic acid and completion of the polymerization. the sebacic acid apparently acts as a chain-extending agent.\nMore recently, U.S. Pat. No. 3,386,967 discloses that a high strength polycaproamide yarn having an excess number of carboxyl end groups over amino end groups may be produced from a high molecular weight polymer prepared by polymerizing e-caprolactam and reacting the polymer with a dibasic carboxylic acid containing at least six carbon atoms. Because polycondensation of the polymer in the presence of a dibasic carboxylic acid is relatively slow, the acid is desirably included at the start of the polymerization reaction. Even so, equilibrium is approached in the polymerization reaction mixture at number average molecular weight not above about 20,000 under usual polymerization conditions, and special measures are required to carry the reaction further. One particularly useful method of accomplishing the required further reaction is to remove volatile by-products of the polymerization such as water by flowing an inert gas across the reaction mixture surface, desirably followed by application of vacuum.\nAlthough the process of U.S. Pat. No. 3,386,967 constitutes a major contribution to this art, we have found that in continuous operation of the process on a commercial basis, particularly at maximum production rates, serious problems have been encountered in melt spinning due to the frequent occurrence of \"nubs\" in the fiber. The term \"nubs\" is conventionally applied and is used herein to mean enlarged sections of filament no more than several filament diameters in length. Nubs may be formed by a foreign, nonorientable substance which interferes with normal fiber stretch in a short section, resulting in an enlargement. Foreign substances which are believed to have contributed to nuts in the present instance include carbonized polymer and gels formed in the polymer. Gels appear to be the chief cause, i.e., the nubs are probably created by non-orientable gel from cross-linked polymer. Thermal degradation of the polymer may be an important causative factor.\nThe reactions in thermal degradation of polyamides containing dicarboxylic acid additives are not entirely understood. It is likely that thermal degradation produces a decomposition product which serves to form cross-links between amide groups and adjacent polymer chains. The decomposition reaction proceeds slowly, finally building up a three-dimensional network of molecules which may be called polymer gel and which eventually reaches the stage where it forms an infusible coating on the walls of the equipment.\nA serious difficulty which arises from the formation of this polymer gel on the interior walls is that from time to time pieces break off and get into the flowing polymer stream where they produce damage to the spinning equipment.\nThe greatest difficulty, however, is caused by polymer gel which has progressed to the three-dimensional structural stage, but which has not yet reached the stage of being infusible. This kind of polymer gel is readily carried with the stream of flowing polymer. Being still molten or at least softened, it passes through the pump and even through the filter medium to show up either as discontinuities or as viscosity differences in the spun filament. When these filaments are later drawn, these defects may cause breaks in the filaments which either cause the whole thread to break or else form nubs which go through to be counted as quality defects in the final yarns."}
{"text": "Polymeric compositions are extensively used as insulation materials for high voltage wire and cable. Olefin homo- and copolymers are the insulation materials of choice for high voltages (5000+ volts) based on desirable electrical and physical properties such as dielectric strength; toughness, such as resistance to cutting and abrasion; ageing characteristics; resilience; and mininal cracking on mechanical stress. Polyethylene, polypropylene and blends thereof polymerized in different densitites as well as copolymers thereof with each other and other comonomers are usually used.\nIt has been noted that such polyolefin-based materials when used as insulation materials in high voltage distribution wires and cable are prone to deterioration under voltage stress in wet environment resulting in electrical failure. The latter deterioration may result from a degradation phenomenon identified as \"water trees\". Under high voltage wet-stress, microscopic channels, i.e. dendritic voids, appear in the insulation. These have a tree-like appearance--hence, the name. Such failure is most disadvantageous. This problem is now aggrevated in that many high voltage cables are buried for greater reliability (decrease of damage from high winds, ice storms, etc.) and for aesthetic reasons. Over extended periods of time, short circits have occurred in such buried cables resulting in loss of service. These cables have to be removed by excavation and replaced, a time consuming and costly operation.\nMany classes of chemical compound additives have been disclosed in the prior art as effective voltage stabilizers, i.e. suppressants for electrical failure, water-treeing and/or electrical-treeing (microscopic dentrites caused by corona arcing). These prior art disclosures include voltage stabilizers based on silicon derivatives, furfuryloxy phosphites and high-molecular weight polyethylene oxide.\nIn the prior art, Kowasaki, U.S. Pat. No. 4,305,849, teaches the use of polyethylene glycols having molecular weights of from about 1,000 to 20,000 as voltage stabilizers.\nAshcraft et al in U.S. Pat. Nos. 4,144,202 and 4,263,158 teach the use of organosilane compounds containing azomethine groups as voltage stabilizers.\nTurbett et al in U.S. Pat. No. 4,376,180 disclose the use of 3-(N-phenylaminopropyl-tridodecyloxysilane) as a voltage stabilizer.\nTurbett in U.S. Pat. No. 4,440,671 discloses the use of a blend of hydrocarbon-substituted diphenyl amine and a high molecular weight polyethylene glycol for this purpose.\nBraus et al, U.S. Pat. No. 4,514,535 disclose the use of tritetrahydrofurfuryloxy phosphite as a voltage stabilizer.\nBeasley et al in U.S. Pat. No. 4,374,224 disclose the use of an organic carboxylic ester having at least one aromatic ring and at least three carboxylic ester groups as a voltage stabilizer.\nU.S. Pat. No. 3,553,348 describes the use of filler minerals such as magnesium silicate, pretreated with alkyl and vinyl alkoxysilanes, as voltage stabilizers.\nIt will be noted from most of these prior art teachings that previous experimenters have relied upon a test based on a microscopic examination of the voltage stressed polymer. The number of \"trees\" and their length were utilized to judge the resistance deterioration of the insulation. These tests were based on the assumption that \"tree\" length and the number of trees could be used to approximate the relative useful life of the insulation. These tests, while rapid, are subject to many variables not quite analagous to the actual service conditions.\nIn contrast, the polymers of this invention and representative compounds of the above mentioned prior art have been tested and compared herein by the use of actual wires coated with the test material compositions. These wires are immersed in a water bath and subjected to high voltage until electrical failure. This test methodology is more closely analagous to the conditions leading to electrical failure of buried distribution cables, and thus provides a meaningful method of illustrating the improved performance characteristics of the instant compounds."}
{"text": "People routinely undergo surgery for treatment of a variety of medical conditions. In many cases, the patients are placed under general anesthesia for the surgical procedure. While under general anesthesia, the patient is unconscious and cannot respond or react to external stimuli or conditions. In some surgical procedures, there is a potential for these external conditions to cause trauma to the patient. For example, during certain orthopaedic surgical procedures, the surgeon must exert considerable pressure on various parts of the body, e.g., to the arms, legs, shoulder, and torso. In addition, the surgeon uses a variety of instruments during the procedure. The pressure on the body and use of surgical instruments create the potential for unintentional trauma to the patient, particularly to the head, face, ears, eyes, and neck area. The patient may receive pressure marks or sores from direct contact with surgical instruments and body positioning techniques. Any pressure applied to the ocular structures, even a short time, can cause damage or blindness to the eye. The body is typically covered at most by a blanket or cloth, which provides minimal protection.\nAnother challenge is the task of rolling the patient over from a supine position to a prone position on the operating table or from a cart onto the operating table. Depending on the procedure, the patient may be intubated, which creates a risk of neck injury during the roll-over process. After the procedure, the patient must again be rolled off the operating table onto a gurney. Still anesthetized, the risk of neck injury is again present if the head is not properly supported and manipulated during patient movement.\nIf an emergency develops while the patient is in the prone position, requiring the patient to be rolled to the supine position, valuable time can be lost trying to properly support the patient without injury to the neck, and without crimping the airway supply tubing and monitoring equipment communicating through the nose and mouth of the patient.\nOne approach found in U.S. Pat. No. 6,490,737 involves a molded helmet, which partially covers the patient's face and head. The helmet is particularly designed for procedures with the patient in the prone position and, in fact, the helmet can be rigidly mounted to the operating table, as shown in FIG. 8 of U.S. Pat. No. 6,490,737. However, in the prone position, the forehead and chin areas are exposed to continuous pressure by the weight of the patient's own head. If not relieved by regular movement of the face to allow blood flow, the pressure can cause localized ischemia to the chin and forehead area. In addition, the helmet leaves the eyes, cheeks, nose, and back of the head exposed and vulnerable to objects smaller than the openings in the helmet. The helmet is molded to a rigid form factor so the head is likely to shift in position relative to the helmet.\nIn the anesthetized state, the patient is unable to respond or react to any of these conditions. The surgeon, anesthesiologist, and surgical staff must be constantly aware of patient safety to avoid unnecessary trauma.\nThe operating room is typically maintained at a low temperature, in part for the comfort and alertness of the surgical team. The patient may experience body heat loss in the low temperature environment. The torso and upper and lower extremities are typically covered by a blanket for warmth. However, a significant amount of heat can be lost passively through the head and neck. A surgical cap can reduce some heat loss, but typically does not cover the face and certainly does not remove the potential for injury."}
{"text": "Without limiting the invention, its background is described in connection with a process for the manufacture of sub-100 nm devices using imprint lithography.\nIn manufacturing, lithography techniques that are used for large-scale production include photolithography and other application oriented lithography techniques, such as electron beam lithography, ion-beam and x-ray lithography, as examples. Imprint lithography is a type of lithography that differs from these techniques. Recent research has shown that imprint lithography techniques can print features that are smaller than 50 nm. As such, imprint lithography has the potential to replace photolithography as the choice for semiconductor manufacturing in the sub-100 nm regime. It can also enable cost effective manufacturing of various kinds of devices, including patterned magnetic media for data storage, micro optical devices, MEMS, biological and chemical devices, X-ray optical devices, etc.\nCurrent research in the area of imprint lithography has revealed a need for devices that can perform orientation alignment motions between a template, which contains the imprint image, and a substrate, which receives the image. Of critical importance is the careful and precise control of the gap between the template and the substrate. To be successful, the gap may need to be controlled within a few nanometers across the imprinting area, while, at the same time, relative lateral motions between the template and the substrate must be eliminated. This absence of relative motion leads is also preferred since it allows for a complete separation of the gap control problem from the overlay alignment problem.\nFor the specific purpose of imprinting, it is necessary to maintain two flat surfaces as close to each other as possible and nearly parallel. This requirement is very stringent as compared to other proximity lithography techniques. Specifically, an average gap of about 100 nm with a variation of less than 50 nm across the imprinted area is required for the imprint process to be successful at sub-100 nm scales. For features that are larger, such as, for example, MEMS or micro optical devices, the requirement is less stringent. Since imprint processes inevitably involve forces between the template and the wafer, it is also desirable to maintain the wafer surface as stationary as possible during imprinting and separation processes overlay alignment is required to accurately align two adjacent layers of a device that includes multiple lithographically fabricated layers. Wafer motion in the x-y plane can cause loss of registration for overlay alignment.\nPrior art references related to orientation and motion control include U.S. Pat. No. 4,098,001, entitled “Remote Center Compliance System;” U.S. Pat. No. 4,202,107, entitled “Remote Axis Admittance System,” both by Paul C. Watson; and U.S. Pat. No. 4,355,469 entitled “Folded Remote Center Compliant Device” by James L. Nevins and Joseph Padavano. These patents relate to fine decoupled orientation stages suitable for aiding insertion and mating maneuvers in robotic machines and docking and assembly equipment. The similarity between these prior art patents and the present invention is in the provision for deformable components that generate rotational motion about a remote center. Such rotational motion is generated, for example, via deformations of three cylindrical components that connect an operator and a subject in parallel.\nThe prior art patents do not, however, disclose designs with the necessary high stiffness to avoid lateral and twisting motions. In fact, such lateral motion is desirable in automated assembly to overcome mis-alignments during the assembly process. Such motion is highly undesirable in imprint lithography since it leads to unwanted overlay errors and could lead to shearing of fabricated structures. Therefore, the kinematic requirements of automated assembly are distinct from the requirements of high precision imprint lithography. The design shown in U.S. Pat. No. 4,355,469 is intended to accommodate larger lateral and rotational error than the designs shown in the first two patents, but this design does not have the capability to constrain undesirable lateral and twisting motions for imprint lithography.\nAnother prior art method is disclosed in U.S. Pat. No. 5,772,905 (the '905 patent) by Stephen Y. Chou, which describes a lithographic method and apparatus for creating ultra-fine (sub-25 nm) patterns in a thin film coated on a substrate in which a mold having at least one protruding feature is pressed into a thin film carried on a substrate. The protruding feature in the mold creates a recess of the thin film. First, the mold is removed from the film. The thin film is then processed such that the thin film in the recess is removed exposing the underlying substrate. Thus, the patterns in the mold are replaced in the thin film, completing the lithography. The patterns in the thin film will be, in subsequent processes, reproduced in the substrate or in another material which is added onto the substrate.\nThe process of the '905 patent involves the use of high pressures and high temperatures to emboss features on a material using micro molding. The use of high temperatures and pressures, however, is undesirable in imprint lithography since they result in unwanted stresses being placed on the device. For example, high temperatures cause variations in the expansion of the template and the substrate. Since the template and the substrate are often made of different materials, expansion creates serious layer-to-layer alignment problems. To avoid differences in expansion, the same material can be used but this limits material choices and increases overall costs of fabrication. Ideally, imprint lithography could be carried out at room temperatures and low pressures.\nMoreover, the '905 patent provides no details relative to the actual apparatus or equipment that would be used to achieve the process. In order to implement any imprint lithography process in a production setting, a carefully designed system must be utilized. Thus, a machine that can provide robust operation in a production setting is required. The '905 patent does not teach, suggest or disclose such a system or a machine.\nAnother issue relates to separation of the template from the substrate following imprinting. Typically, due to the nearly uniform contact area at the template-to-substrate interface, a large separation force is needed to pull the layers apart. Such force, however, could lead to shearing and/or destruction of the features imprinted on the substrate, resulting in decreased yields.\nIn short, currently available orientation and overlay alignment methods are unsuitable for use with imprint lithography. A coupling between desirable orientation alignment and undesirable lateral motions can lead to repeated costly overlay alignment errors whenever orientation adjustments are required prior to printing of a field (a field could be for example a 1″ by 1″ region of an 8″ wafer).\nFurther development of precise stages for robust implementation of imprint lithography is required for large-scale imprint lithography manufacturing. As such, a need exists for an improved imprint lithography process. A way of using imprint lithography as a fabrication technique without high pressures and high temperatures would provide numerous advantages."}
{"text": "An important goal in animal husbandry is to develop biologically active agents which can increase the quantity and improve the quality of meat obtained from livestock animals.\nxe2x80x9cIncreasing the quantityxe2x80x9d of food obtained from livestock animals refers to, inter alia, promoting the growth of livestock animals, increasing the efficiency of feed utilized in raising livestock animals and/or enhancing the production of lean body mass in livestock animals. Biologically active agents with these attributes are commonly referred to as xe2x80x9canabolic agentsxe2x80x9d.\nxe2x80x9cImproving the qualityxe2x80x9d of food obtained from livestock animals refers to, inter alia, reducing the quantity of subcutaneous fat in meat and, in poultry, the size of the abdominal xe2x80x9cfat padxe2x80x9d. Subcutaneous fat can cause elevated cholesterol and triglyceride levels in individuals who consume large quantities of meat, has minimal nutritional value and decreases the overall yield of meat. Therefore, the reduction or elimination of this type of fat from meat is desirable. On the other hand, intramuscular fat, commonly referred to as xe2x80x9cmarblingxe2x80x9d, contributes positively to the flavor of meat and maintains a high Quality Grade. Marbling is therefore considered a desirable quality. Biologically active agents which are lipolytic can reduce subcutaneous fat while retaining the intramuscular fat.\nCertain publications have appeared generally disclosing arylpropanolamines such as U.S. Pat. No. 5,013,761 and WO 97/10825. There continues to be a need for the development of agents which are anabolic and lipolytic to improve the economics of meat production by increasing the yield and improving the quality of meat obtained from livestock animals. Further increases in profitability could be achieved by the development of long-lasting anabolic/lipolytic agents which, because they require less frequent dosing, are more convenient and economical to use.\nIt has now been found that indazolyloxy propanolamimines such as those represented by Structural Formula (I) and (III) below are both anabolic and lipolytic when administered to livestock (Example 7 and Table 1). It has also been found that the anabolic and lipolytic effects of these compounds are longer lasting than other aryloxy propanolamines. Specifically, the percent decrease in serum urea nitrogen level (SUN) after forty-eight hours and the increase in non-esterified fatty acid levels (NEFA) after twenty-four hours in cattle treated with the indazolyloxy propanolamines of the present invention have been found to be significantly greater than for other aryloxy propanolamines (see Example 8 and Table 2). The percent decrease in SUN is indicative of anabolic activity and the increase in NEFA is indicative of lipolytic effects. Based on these results, novel compounds and novel methods of improving meat production from livestock animals are disclosed herein.\nOne embodiment of the present invention is a compound represented by Structural Formula (I): \nand physiologically acceptable salts thereof; where R1 and R2 are independently xe2x80x94H or a C1-C4 straight chained or branched alkyl group;\nRing A, Ring B and Ring C are independently substituted or unsubstituted, provided, however, that Ring C is not substituted in the para position with a group represented by Structural Formula (II): \nR3 and R4 are independently xe2x80x94H, a straight or branched C1-C4 alkyl or, taken together with the nitrogen atom to which they are bonded, form a non-aromatic heterocyclic ring, and Ring D is substituted with zero, one, two or three additional substituents.\nAnother embodiment of the present invention is a compound represented by Structural Formula (III): \nand physiologically acceptable salts thereof; where\nR1 and R2 are independently xe2x80x94H, a C1-C4 straight chained or a branched alkyl group;\nRing A, Ring B, Ring C and Ring E are independently substituted or unsubstituted, provided, however, that Ring E is not substituted in the position meta to xe2x80x94Nxe2x80x94 and ortho to the carbon bonded to oxygen with xe2x80x94CONR3R4;\nR3 and R4 are independently xe2x80x94H or a straight or branched chain C1-C4 alkyl group or, taken together with the nitrogen atom to which they are bonded, form a non-aromatic heterocyclic ring.\nAnother embodiment of the present invention is a method of increasing the quantity and improving the quality of meat obtained from a livestock animal. The method comprises administering to the animal an effective amount of one or more compounds represented by Structural Formula (I). Another embodiment of the present invention is a method of increasing the quantity and improving quality of meat obtained from a livestock animal. The method comprises administering to the animal an effective amount of one or more compounds represented by Structural Formula (III).\nThe indazolyloxy propanolamines of the present invention have both anabolic and lipolytic activity and can therefore be administered to livestock to increase the output of meat and to decrease its fat content. Moreover, the duration of the anabolic and lipolytic activity of these indazolyloxy propanolamines is longer in vivo than other structurally related aryloxy propanolamines. Consequently, indazolyloxy propanolamines are more convenient and economical to use than other aryloxy propanolamines because indazolyloxy propanolamines need to be administered less frequently and/or at lower doses. The compounds of the present invention are intended for the treatment of healthy animals.\nThe present invention is directed to a compound represented by Structural Formula (I) and (III). Also included is a method of improving livestock production by administering one or more compounds of the present invention to the livestock.\nIn a preferred embodiment, the compound is represented by Structural Formula (IV): \nRing A, Ring B, Ring C, R1 and R2 are as described above for Structural Formula (I). Ring E is substituted or unsubstituted.\nX is xe2x80x94CHxe2x80x94 or xe2x80x94Nxe2x80x94.\nIn Structural Formula (IV), R1 and R2 are preferably methyl and Rings A-C have no further substitution.\nIn a preferred embodiment, the compound of the present invention is represented by Structural Formula (V): \nX is xe2x80x94Nxe2x80x94 or xe2x80x94CHxe2x80x94.\nRing E is substituted or unsubstituted. Examples of suitable substituents for Ring E include halo, xe2x80x94CN, xe2x80x94OR5, C1-C4 alkyl, C1-C4 haloalkyl, xe2x80x94CO2R5, xe2x80x94CONR6R7, xe2x80x94CONH(C1-C4 alkyl), xe2x80x94SR5, xe2x80x94CSNR6R7, xe2x80x94CSNR6R7, xe2x80x94SO2R5, xe2x80x94SO2NR6R7, xe2x80x94SOR5, xe2x80x94NR6R7. Preferred substituents are xe2x80x94CN, xe2x80x94CONH2, xe2x80x94SO2CH3, xe2x80x94SO2NH2 or are represented by Structural Formulas (VI) or (VII): \nxe2x80x94SR2CH3 is a more preferred substituent.\nR5 is xe2x80x94H, C1-C4 alkyl or aryl.\nR6 and R7 are independently xe2x80x94H, C1-C4 alkyl, aryl, xe2x80x94(CH2)naryl, or combine with the nitrogen atom to which each is bound to form morpholinyl, piperidinyl, pyrrolidinyl, or piperazinyl.\nn is 0, 1, 2, or 3.\nIn another preferred embodiment, the compound of the present invention is represented by Structural Formula (IV) or (V), provided, however, that Ring E is not substituted in the position meta to xe2x80x94Xxe2x80x94 and ortho to the carbon bonded to oxygen with xe2x80x94CONR3R4. R3 and R4 are independently xe2x80x94H or a straight or branched chain C1-C4 alkyl group or, taken together with the nitrogen atom to which they are bonded, form a non-aromatic heterocyclic ring.\nAlso included in the present invention are Compounds 1-5, shown in Tables 1 and 2.\nPhysiologically acceptable salts of the compounds disclosed herein, including the compounds represented by Structural Formulas (I),(III), (IV), (V) and Compounds 1-5, shown in Tables 1 and 2, are included. Salts can be formed from those compounds which comprise acidic functional groups by reacting with a suitable base. Such salts include those derived from inorganic bases such as ammonium and alkali and alkaline earth metal hydroxides, carbonates, bicarbonates, and the like, as well as salts derived from basic organic amines such as aliphatic and aromatic amines, aliphatic diamines, hydroxy alkamines, and the like. Such bases useful in preparing the salts of this invention thus include ammonium hydroxide, potassium carbonate, sodium bicarbonate, calcium hydroxide, methylamine, diethylamine, ethylenediamine, cyclohexylamine, ethanolamine and the like.\nBecause of the amine moiety, salts of the compounds disclosed herein can also be prepared by reacting with a suitable acid. Acids commonly employed to form such salts include inorganic acids such as hydrochloric, hydrobromic, hydroiodic, sulfuric and phosphoric acid, as well as organic acids such as para-toluenesulfonic, methanesulfonic, oxalic, para-bromophenylsulfonic, carbonic, succinic, citric, benzoic, acetic acid, and related inorganic and organic acids. Such pharmaceutically acceptable salts thus include sulfate, pyrosulfate, bisulfate, sulfite, bisulfite, phosphate, mono-hydrogenphosphate, dihydrogenphosphate, metaphosphate, pyrophosphate, chloride, bromide, iodide, acetate, propionate, decanoate, caprylate, acrylate, formate, isobutyrate, heptanoate, propiolate, oxalate, malonate, succinate, suberate, sebacate, fumarate, 2-butyne-1,4 dioate, 3-hexyne-2, 5-dioate, benzoate, chlorobenzoate, hydroxybenzoate, methoxybenzoate, phthalate, xylenesulfonate, phenylacetate, phenylpropionate, phenylbutyrate, citrate, lactate, hippurate, xcex2-hydroxybutyrate, glycolate, maleate, tartrate, methanesulfonate, propanesulfonate, naphthalene-1-sulfonate, naphthalene-2-sulfonate, mandelate and the like salts.\nWhen substituted, Rings A-E can have one, two or three substituents in addition to those shown in Structural Formulas (I)-(V). Suitable substituents are those which do not significantly decrease the anabolic and lipolytic properties of the compound. Examples of suitable substituents include halogens, hydroxy, xe2x80x94ORxe2x80x2, xe2x80x94SRxe2x80x2, xe2x80x94S(O)Rxe2x80x2, xe2x80x94S(O)2Rxe2x80x2, xe2x80x94COORxe2x80x2, xe2x80x94C(O)Rxe2x80x2, xe2x80x94CN, xe2x80x94NO2, xe2x80x94OCONRxe2x80x2Rxe2x80x3, xe2x80x94OCONHRxe2x80x2, xe2x80x94NHCOORxe2x80x2, xe2x80x94NRxe2x80x3COORxe2x80x2, xe2x80x94NHRxe2x80x2, xe2x80x94NRxe2x80x2Rxe2x80x3, xe2x80x94CN, C1-C4 alkyl, C1-C4 haloalkyl, xe2x80x94CONRxe2x80x2Rxe2x80x3, xe2x80x94CONHRxe2x80x2, xe2x80x94CSNRxe2x80x2Rxe2x80x3, xe2x80x94CSNRxe2x80x2Rxe2x80x3, xe2x80x94SO2NRxe2x80x2Rxe2x80x3. Rxe2x80x2 and Rxe2x80x3 are independently xe2x80x94H, C1-C4 alkyl or aryl. In addition, when Rxe2x80x2 and Rxe2x80x3 are bonded to the same nitrogen atom (e.g, xe2x80x94NRxe2x80x2Rxe2x80x3, xe2x80x94CONRxe2x80x2Rxe2x80x3 or xe2x80x94OCONRxe2x80x2Rxe2x80x3), then R and Rxe2x80x2, taken together with the nitrogen atom, can form a non-aromatic heterocyclic ring.\nAryl groups include carbocyclic aromatic groups such as phenyl, 1-naphthyl, 2-naphthyl, 1-anthracyl and 2-anthracyl, and heteroaryl groups such as N-imidazolyl, 2-imidazolyl, 2-thienyl, 3-thienyl, 2-furanyl, 3-furanyl, 2-pyridyl, 3-pyridyl, 4-pyridyl, 2-pyrimidyl, 4-pyrimidyl, 2-pyranyl, 3-pyranyl, 3-pyrazolyl, 4-pyrazolyl, 5-pyrazolyl, 2-pyrazinyl, 2-thiazolyl, 4-thiazolyl, 5-thiazolyl, 2-oxazolyl, 4-oxazolyl and 5-oxazolyl.\nAryl groups also include fused polycyclic aromatic ring systems in which a carbocyclic aromatic ring or heteroaryl ring is fused to one or more other heteroaryl rings. Examples include 1-benzimidazolinyl, 2-benzimidazolonyl, 1-benzimidthioazolinyl, 2-benzimidthioazolonyl, 1-carbazolyl, 2-carbazolyl, 3-carbazolyl, 4-carbazolyl, 3-indazolyl, 4-indazolyl, 5-indazolyl, 6-indazolyl, 2-benzothienyl, 3-benzothienyl, 2-benzofuranyl, 3-benzofuranyl, 2-indolyl, 3-indolyl, 2-quinolinyl, 3-quinolinyl, 2-benzothiazolyl, 2-benzooxazolyl, 2-benzimidazolyl, 2-quinolinyl, 3-quinolinyl, 1-isoquinolinyl, 3-quinolinyl, 1-isoindolyl and 3-isoindolyl. Also included within the scope of the term aryl group, as it is used herein, is a group in which one or more carbocyclic aromatic rings and/or heteroaromatic rings are fused to a cycloalkyl or non-aromatic heterocyclic ring.\nNon-aromatic heterocyclic rings are non-aromatic carbocyclic rings which include one, two or three heteroatoms selected from nitrogen, oxygen and sulfur in the ring that will afford a stable structure. The ring can be five, six, seven or eight-membered. Examples include 2-tetrahydrofuranyl, 3-tetrahydrofuranyl, 2-tetrahyrothiophenyl, 3-tetrahyrothiophenyl, 2-morpholino, 3-morpholino, 4-morpholino, 2-thiomorpholino, 3-thiomorpholino, 4-thiomorpholino, 1-pyrrolidinyl, 2-pyrrolidinyl, 3-pyrrolidinyl, 1-piperazinyl, 2-piperazinyl, 1-piperidinyl, 2-piperidinyl, 3-piperidinyl, 4-piperidinyl and 4-thiazolidinyl.\nAs used herein, aliphatic groups include straight chained, branched or cyclic C1-C20 hydrocarbons which are completely saturated or which contain one, two or three units of unsaturation.\nSuitable substituents on an aliphatic group, aryl group (carbocyclic and heteroaryl), non-aromatic heterocyclic ring or benzyl group are those which do not significantly reduce the anabolic effects or alter the lipolytic effects of the compound. Examples include xe2x80x94OH, halogen (xe2x80x94Br, xe2x80x94Cl, xe2x80x94I and xe2x80x94F), xe2x80x94OR, xe2x80x94Oxe2x80x94COR, xe2x80x94CN, xe2x80x94NO2, xe2x80x94COOH, xe2x80x94NH2, xe2x80x94NHR, xe2x80x94NR2, xe2x80x94COOR, xe2x80x94COR, xe2x80x94CHO, xe2x80x94CONH2, xe2x80x94CONHR, xe2x80x94CONR2, xe2x80x94SH, xe2x80x94SR and xe2x80x94NHxe2x80x94C(xe2x95x90NH)xe2x80x94NH2. R is C1-C6 alkyl, benzyl, or phenyl.\nA substituted non-aromatic heterocyclic ring can also have xe2x95x90O, xe2x95x90S, xe2x95x90NH or xe2x95x90NR where R is as defined above as a substituent. A substituted aliphatic, substituted aromatic, substituted non-aromatic heterocyclic ring or substituted benzyl group can have one, two or three substituents.\nIn the structural formulas depicted herein, the bond by which a chemical group or moiety is connected to the remainder of the molecule or compound is indicated by the following symbol:\n\nFor example, the corresponding symbol in Structural Formula (II) indicates the bond by which the oxygen of the pyridyloxy group is connected to Ring C of Structural Formula (I).\nThe present invention includes solvates of the compounds of Structural Formula I and the physiologically acceptable salts thereof. A particular compound of the present invention or a physiologically acceptable salt thereof may form solvates with water or common organic solvents. Such solvates are included within the scope of compounds of the present invention.\nIn addition, it will be appreciated that diastereomers exist for the compounds of Structural Formula I and, depending on the substituents, further diastereomers may exist. The compounds of the present invention include mixtures of two or more diastereomers as well as each individual stereoisomer.\nIt will also be appreciated that some of the heterocycles may exist in tautomeric forms. All such forms are included within the scope of the present invention.\nAs used herein, xe2x80x9ccrystallizingxe2x80x9d refers to providing a solvent or solvent mixture in which the indazolyloxy propanolamine compound is highly soluble and in which an ammonium salt(s) thereof is insoluble or only slightly soluble. The compound is dissolved in the solvent or solvent mixture and then converted to the ammonium salt by the addition of at least one equivalent of the appropriate acid, after which the ammonium salt precipitates. To minimize contamination of the precipitated product, between about 1.0 and about 1.1 equivalents of acid are preferably used. Impurities present in the indazolyloxy propanolamine are preferably highly soluble in the solvent or solvent mixture, resulting in a precipitated salt which is purified relative to the free base prior to precipitation. More preferably, the precipitated compound is crystalline.\nPurification by acidic (or basic) extraction refers to dissolving a compound with a basic functional group such as an amine (or a compound with an acidic functional group, such as a carboxylic acid) in aqueous or alcoholic acid (or aqueous or alcoholic base, in the case of a compound with an acidic functional group). The aqueous solution can then be washed with organic solvents that are not miscible with water to remove organic impurities. The pH of the solution is then adjusted to the isoelectric point of the compound, thereby precipitating the compound or allowing its extraction into an organic solvent.\nPurification by precipitation at the isoelectric point refers to adjusting the pH of an aqueous solution of a compound with both an acidic and basic functional group (e.g., an amino acid) to its isoelectric point, thereby causing the compound to precipitate from solution. xe2x80x9cIsoelectric pHxe2x80x9d is the pH at which the compound is electrically neutral and therefore least soluble in aqueous solution. Preferably, impurities present in the compound are highly soluble at the isoelectric point. As a result, the precipitated compound is purified relative to the compound prior to precipitation. More preferably, the precipitated compound is crystalline. Optionally, a compound with both an acidic and basic functional group can be purified by both by acidic (or basic) extraction and by precipitation at its isoelectric point.\nLivestock animals are animals raised for food production. Ruminants or xe2x80x9ccud-chewingxe2x80x9d animals such as cows, bulls, heifers, steers, sheep, buffalo, bison, goats and antelopes are examples of livestock. Other examples of livestock include pigs and avians (poultry) such as chickens, ducks, turkeys and geese. Yet other examples of livestock include fish, shellfish and crustaceans raised in aquaculture. Also included are exotic animals used in food production such as alligators, water buffalo and ratites (e.g., emu, rheas or ostriches). The method of the present invention is preferably used with avians.\nAn xe2x80x9ceffective amountxe2x80x9d of a compound of the present invention is the quantity which, when administered to a livestock animal, increases the quantity of meat and/or quality of meat obtained from the animal.\nIncreasing the quantity of meat obtained refers to promoting a greater amount of growth in the animal with a treatment compared with the absence of the treatment. Alternatively, increasing the quantity of meat obtained refers to promoting formation of lean body mass. The formation of lean body mass is promoted, for example, when there is a higher ratio of muscle to fat as a result of a treatment than in the absence of the treatment. Alternatively, increasing the quantity of meat obtained refers to improving the efficiency of utilization of food. Food utilization is more efficient when there is a greater body weight gain per a given amount of feed consumed by an animal as a result of a treatment than in its absence. Thus, increasing the quantity of meat obtained from a livestock animal generally results in an improvement in the economics, e.g., in increase in the profitability of producing meat.\nIncreasing the quality of meat refers to an improvement in carcass quality of the animal. Improved carcass quality refers, for example, to the formation of less fatty tissue (subcutaneous fat), to a decreased size of the fat pad in poultry and/or to greater leanness (improved yield). Thus, improved carcass quality generally results in meat that is more healthy to consume, e.g., is less likely to cause elevated cholesterol and/or triglyceride levels. Improving the quality of meat can also improve the economics and increase the profitability of producing meat because high quality grade meat can command higher selling prices at market.\nThe effective amount to be administered will vary somewhat depending upon the particular animal species being treated and the particular active ingredient employed, but generally will be from about 0.5 to about 1000 parts per million (ppm: milligrams compound per kilogram food) of total daily feed intake. Such amount will provide a dosage of about 0.02 to about 50 mg/kg. A preferred embodiment employs about 0.5 to about 200 ppm, and more preferably from about 1 to about 40 ppm. For example, when practicing the method in animals such as poultry, the compound will be added to the daily feed ration at about 2 to 100 parts per million of the daily feed ration.\nThe method of the present invention is preferably practiced by orally administering an effective amount of a compound of the present invention to a livestock animal. Other routes of administration can be employed, for instance intranasal (e.g., by intranasal misting device), in ovo or subcutaneous, intramuscular or intravenous injection; however, such routes are less practical.\nFor oral administration, a compound of the present invention is preferably admixed with suitable carriers or diluents commonly employed in animal husbandry. Animal feedstuffs comprising a compound of the present invention are provided as a further embodiment of this invention. Typical carriers and diluents commonly employed in such feedstuffs include corn meal, corncob grits, soybean meal, alfalfa meal, rice hulls, soybean mill run, cottonseed oil meal, bone meal, ground corn, corncob meal, wheat middlings, limestone, dicalcium phosphate, sodium chloride, urea, distillers dried grain, vitamin and/or mineral mixes, cane molasses or other liquid carriers and the like. Such carriers promote a uniform distribution of the active ingredient, and more typically about 20 to about 98 percent by weight.\nWhile the preferred method for orally administering the compounds of the present invention is via the daily feed rations, the compounds can be incorporated into salt blocks and mineral licks, as well as being added directly to link tank formulations or drinking water for convenient oral consumption. The compounds can additionally be formulated with polymorphous materials, waxes and the like for long-term controlled release, and administered to animals as a bolus or tablet only as needed to maintain the desired daily payout of active ingredient. Compounds can also be administered orally by gavage treatment and/or applied transdermally.\nFor parenteral administration, the compounds of the present invention can be admixed with conventional carriers such as water, propylene glycol, polyethylene glycols, n-methyl pyrrolidone, glycerol formal, corn oil, sesame oil, calcium stearate, polymeric materials and the like. Such formulations can be molded into pellets and administered as an injection or as a slow-release subcutaneous implant, sustained rumen delivery device or intranasal device. Such administrations can be made as often as needed to ensure the proper dosing of active ingredient to obtain the desired rate of growth promotion and improvement in leanness and feed efficiency.\nThe compounds of the present invention can be prepared by procedures disclosed in WO 97/10825 to Bell et al., WO 98/09625 to Crowell et al., U.S. Pat. Nos. 5,808,080 and 6,046,227. The entire teachings of these references are incorporated herein by reference. A reaction scheme for preparing these compounds is shown below: \nX and Ring D in the Scheme are as described above.\nThe amination of the epoxides in the Scheme is carried out under conditions known in the art for this type of reaction. For example, the epoxide may be combined with the amine in an alcohol, preferably ethanol at room temperature, to the reflux temperature of the reaction mixture. For example, the reaction is carried out under conditions generally described in Atkins et al., Tetrahedron Letters 27:2451 (1986) the entire teachings of which are incorporated herein by reference. An example of specific conditions for reacting an epoxide with an amine is provided in Example 6.\nSubstituents which interfere with the reaction shown in the Scheme can be present, provided that they are first converted to a protected form. Suitable protecting groups are known to those skilled in the art and are disclosed in Green and Wuts, xe2x80x9cProtecting Groups in Organic Synthesisxe2x80x9d, John Wiley and Sons, 1991, the teachings of which are incorporated herein by reference.\nAn alternative method of preparing the compounds of the present invention involves deprotecting a cyclic sulfate-containing compound to reveal hydroxy substituents, and comprises combining the cyclic sulfate-containing compound with a trialkylsilyl halide in a solvent for a time sufficient to deprotect the hydroxyl group.\nSpecifically, the process comprises reacting a compound of the formula (VIII): \nwith thionyl halide in a solvent for a time sufficient to yield a sulfite compound of the formula (IX): \nwhere Ring A and Ring B are as described above.\nSuitable thionyl halides include thionyl bromide and thionyl chloride. Thionyl chloride is preferred. Typically, 2 equivalents of thionyl halide are used per mole of compound (VIII).\nAny solvent can be used in this step, so long as it does not interfere with the reaction. THF is preferably used. The reaction is typically conducted until product is obtained, generally for about 120 to 240, preferably 180, minutes. The reaction can be conducted at any temperature, but generally is conducted at a temperature less than about 0xc2x0 C., preferably from about xe2x88x9210 to 0xc2x0 C., more preferably from about xe2x88x9210 to xe2x88x925xc2x0 C., most preferably from about xe2x88x929 to xe2x88x928xc2x0 C.\nIn a preferred embodiment, the NH group in the 1H-indazole ring is first N-protected prior to conversion of the compound to the corresponding sulfate. Selection of the protecting group is preferably made such that it can be removed under the same conditions as for removing the sulfate group to reveal a hydroxy group.\nConditions for N-protecting the nitrogen atoms depend on the particular protecting group chosen. Suitable reaction conditions are described in xe2x80x9cProtective Groups in Organic Synthesis,xe2x80x9d Peter G. M. Wuts (Editor), Theodora W. Greene, 3rd ed. (April 1999), Vch Pub.\nThereafter, compound (IX) can be converted into corresponding sulfate compound of the formula (X): \nwhere Ring A and Ring B are as described above.\nSuitably, the conversion comprises combining sulfite compound (IX) with a catalytic amount of a ruthenium compound and an oxidizing agent in solvent for a time sufficient to oxidize said sulfite compound (IX) to a sulfate compound (X).\nSuitable ruthenium compounds include ruthenium chloride or ruthenium oxide. Generally, about 0.001 to 0.25 equivalents of ruthenium compound are used per mole of sulfite compound (IX).\nSuitable oxidizing agents include sodium periodate, sodium hypochlorite, sodium bromate, calcium hypochloride, sodium chlorate or ozone. Generally, about greater than about 2.0, preferably 2.5, equivalents of oxidizing agent are used per mole of sulfite compound (IX).\nThe conversion is typically conducted in any solvent that would not interfere with the reaction, such as CCl4, CHCl3, CH3CN, or water or mixtures thereof. The conversion is typically conducted for about 30 to 120, preferably 60, minutes. Preferably, the conversion is conducted in a mixture of CHCl3, CH3CN and water. The conversion is typically conducted at a temperature of from about xe2x88x9210 to 25xc2x0 C.\nThereafter, the sulfate compound (X) can be reacted with a primary amine of the formula (XI): \nwhere X, Ring C and Ring E are as described above, in solvent for a time sufficient to yield a compound (XII): \nTypically, about 0.9 to 2 equivalents of amine are used per mole of sulfate compound (X). Preferably, about 0.9 to 1.1 eq. are used. Most preferably, about 1.0 eq. (a stoichiometric amount) of amine is used.\nTypically, any solvent can be used in this step, as long as it does not interfere with the reaction. Preferably, the reaction is conducted in CH3CN. The reaction is typically conducted for about 60-180, preferably 120, minutes at a temperature of from about 78 to 85xc2x0 C.\nFinally, the compound (XII) can be combined with trialkylsilyl halide in solvent for a time sufficient to yield compound (IV): \nwhere R1, R2, X, Ring A, Ring B, Ring C, and Ring E are as described above.\nSuitable trialkylsilyl halides include trialkylsilyl iodide, trialkylsilyl bromide or trialkylsilyl chloride. Preferably, the trialkylsilyl halide is trimethylsilyl iodide.\nThe reaction may be carried out in any solvent, so long as the solvent does not interfere with the reaction. Preferably, the solvent is an aprotic solvent such as carbon tetrachloride, acetonitrile, dimethyl sulfoxide, dimethylformamide, sulfolane, tetrahydrofuran, diethyl ether, methyl-t-butyl ether, 1,2-dimethoxy-ethane, dioxane, chloroform, methylene chloride, toluene or acetone or mixtures thereof.\nThe reaction is typically conducted for 5-25, preferably 15, minutes at a temperature of less than 0xc2x0 C., preferably xe2x88x9210 to 0xc2x0 C., most preferably xe2x88x929 to xe2x88x928xc2x0 C.\nA method of preparing the compounds of the present invention involves formation of salts of novel indazolyloxy propanolamines in which the amine group is substituted with an alkyl 2-[4-(2-yl-2-methylpropyl)phenoxy]pyridine-carboxylate group (hereinafter xe2x80x9cindazolyloxy propanolamine estersxe2x80x9d). These indazolyloxy propanolamine esters can be readily crystallized in high yield and in high purity, substantially free (typically less than 1 ppm) of epoxide precursors. In addition, the corresponding carboxylic acids of indazolyloxy propanolamine esters (hereinafter xe2x80x9cindazolyloxy propanolamine carboxylic acidsxe2x80x9d) can be purified by precipitation at their isoelectric point in high yield and high purity. Both indazolyloxy propanolamine esters and carboxylic acids can be prepared in high yield from readily accessible starting materials.\nThis method can be used to prepare an indazolyloxy propanolamine ester represented by Structural Formula (XIII): \nand ammonium salts thereof.\nR8 is a C1 to C6 straight or branched chain alkyl group or a C7 to C9 substituted or unsubstituted aralkyl group. Preferably, R8 is methyl or ethyl.\nThe method comprises the step of reacting an epoxide starting material with an amine starting material. The epoxide starting material is represented by Structural Formula (XIV): \nand the amine starting material is represented by Structural Formula (XV):. \nThis method of can also be used to prepare an indazolyloxy propanolamine carboxylic acid or a carboxylate salt thereof from an indazolyloxy propanolamine ester represented by Structural Formula (XIII). The indazolyloxy propanolamine carboxylic acid is represented by Structural Formula (XVI): \nThe method comprises the step of hydrolyzing the xe2x80x94COOR8 group of the indazolyloxy propanolamine ester.\nBecause of the amine moiety present in the indazolyloxy propanolamines esters disclosed herein, ammonium salts of these compounds can be prepared by reacting with a suitable acid. Acids commonly employed to form such salts include inorganic acids such as hydrochloric, hydrobromic, hydroiodic, sulfuric and phosphoric acid, as well as organic acids such as para-toluenesulfonic, methanesulfonic, oxalic, para-bromophenylsulfonic, carbonic, succinic, citric, benzoic, acetic acid, and related inorganic and organic acids. Such pharmaceutically acceptable salts thus include sulfate, pyrosulfate, bisulfate, sulfite, bisulfite, phosphate, mono-hydrogenphosphate, dihydrogenphosphate, metaphosphate, pyrophosphate, chloride, bromide, iodide, acetate, propionate, decanoate, caprylate, formate, isobutyrate, heptanoate, propiolate, oxalate, malonate, succinate, suberate, sebacate, fumarate, maleate, 2-butyne-1,4 dioate, 3-hexyne-2, 5-dioate, benzoate, chlorobenzoate, hydroxybenzoate, methoxybenzoate, phthalate, xylenesulfonate, phenylacetate, phenylpropionate, phenylbutyrate, citrate, lactate, hippurate, xcex2-hydroxybutyrate, glycolate, maleate, tartrate, methanesulfonate, propanesulfonate, naphthalene-1-sulfonate, naphthalene-2-sulfonate and mandelate salts.\nCarboxylate salts can be formed from the indazolyloxy propanolamines acids disclosed herein which have a carboxylic acid functional group by reacting with a suitable base. Such salts include those derived from inorganic bases such as ammonium and alkali and alkaline earth metal hydroxides, carbonates, bicarbonates, and the like, as well as salts derived from basic organic amines such as aliphatic and aromatic amines, aliphatic diamines, hydroxy alkylamines, and the like. Such bases useful in preparing the salts of this invention thus include ammonium hydroxide, potassium carbonate, sodium bicarbonate, calcium hydroxide, sodium hydroxide, lithium hydroxide, potassium hydroxide, methylamine, diethylamine, ethylenediamine, cyclohexylamine, ethanolamine and the like.\nPreferred carboxylate salts of the compound represented by Structural Formula (XVI) include alkali metal carboxylate salts such as the sodium carboxylate salt, the potassium carboxylate salt and the lithium carboxylate salt.\nThe preparation of the indazolyloxy propanolamine ester represented by Structural Formula (XIII) from the epoxide and amine represented by Structural Formulas (XIV) and (XV), respectively, is generally carried out in a solvent at room temperature to the reflux temperature of the reaction mixture. Temperatures of 40-140xc2x0 C. are generally preferred. The solvents that may be used in the reaction include: alcoholic solvents, such as methanol, ethanol or isopropanol, (with the preferred solvent corresponding to the particular ester being used in the reaction to prevent transesterification); aromatic solvents, such as benzene, toluene, xylenes, chlorobenzene, dichlorobenzene, other haloaromatics, nitrobenzene, benzonitrile, or trifluoromethylbenzene; or dipolar aprotic solvents, such as dimethylsulfoxide (DMSO), N,N-dimethylacetamide (DMAC) or N,N-dimethylformamide; or other solvents where the reagents are soluble and the temperature can be elevated to the range previously described. Equimolar amounts of epoxide and amine can be used. Alternatively, up to about a five fold excess of one reagent is used. Preferably, however, between about 1.1 to about 2.0 equivalents of amine relative to epoxide is preferred.\nIndazolyloxy propanolamine esters, prepared as described above, can be purified by any suitable means, including by converting to a suitable ammonium salt and crystallizing, as described above.\nAmine starting materials can be prepared as generally described in Examples 3 and 4. A phenoxide salt of the alkylamino phenol represented by Structural Formula (XVII) is reacted with a 2-halo pyridine ester represented by Structural Formula (XVIII): \nR8 in Structural Formula (XVIII) is as described for Structural Formula (XIII). X is a halo group. The coupling reaction and the preparation of the compound represented by Structural Formula (XVII) can be carried out according to procedures described in the aforementioned WO 97/10825 to Bell et al., WO 98/09625 to Crowell et al., U.S. Pat. Nos. 5,808,080 and 6,046,227.\nGenerally, the reaction may be carried out by mixing the alkylamino phenol with a base in the presence of the 2-halopyridine ester. Equimolar amounts of the starting materials are preferably used. However, molar excesses up to about five or ten fold of one starting material relative to the other can be used. The coupling is performed by mixing the amino phenol with a base in a suitable solvent or solvent system. The reaction can be carried out at temperatures as low as room temperature, but is preferably carried out by heating the mixture at reflux while azeotropically removing water formed during the deprotonation step. The 2-halopyridine ester is then added and the reaction continued until the reaction is complete.\nSuitable solvents include dipolar aprotic, ethereal, and aromatic solvents, as well as combinations thereof. Dipolar aprotic solvents include solvents such as DMSO, N,N-dimethylacetamide, N-methylpyrrolidinone, 1,3-dimethyl-2-imidazolindinone (DMI), 1,3-dimethyl-3,4,5,6-tetrahydro-2(1H)-pyrimidinone (DMPU) with a lower boiling solvent to azeotropically remove water such as benzene, toluene, isooctane, xylenes or other solvents capable of forming binary azeotropes with water but which are inert under the reaction conditions. Other suitable compounds for azeotropic removal of water may be found in Advances In Chemistry Series 116xe2x80x94Azeotropic Data III, American Chemical Society: Washington D.C., 1973. Ethereal solvents include tetrahydrofuran, dioxane and 1,2-dimethoxyethane. Aromatic solvents include benzene, toluene, chlorobenzene, anisole and 1,2-dichlorobenzene. Preferred is the use of an aromatic solvent such as chlorobenzene containing 0.1-10 equivalents of a dipolar aprotic solvent such as N,N-dimethylacetamide. Suitable bases include alkali metal alkoxides, such as alkali metal methoxides, ethoxides and tert-butoxides (preferably corresponding to the alkyl group of the ester to prevent transesterification), alkali metal hydroxides, such as lithium hydroxide, sodium hydroxide and potassium hydroxide, or alkali metal carbonates, such as Na2CO3 or K2CO3. Although excesses of base up to at least about ten fold can be used, preferred is the use of stoichiometric quantities of hydroxide or alkoxide bases or 1.5-5 equivalents of K2CO3.\nThe preparation of the epoxide starting material represented by Structural Formula (XIV) can be carried out by reacting 4-hydroxyindazole with (2S)-(+)-glycidyl 3-nitrobenzenesulfonate in an inert solvent (e.g., acetone, methyl ethyl ketone, methyl isobutylketone, dimethyl sulfoxide, N,N-dimethylacetamide (DMAC) or DMF) and in the presence of a base. Although about equimolar amounts of the starting materials are preferred, molar excess up to about five to about ten-fold of one starting material relative to the other can also be used. Suitable bases include non-nucleophilic bases such as potassium carbonate, sodium carbonate and alkali metal alkoxides; potassium carbonate is preferred. Although an excess of base can be used, between about 1.05 and about 1.50 equivalents of K2CO3 are preferred. The reaction is carried out at temperatures ranging from about ambient temperature to about 40xc2x0 C., preferably about 30xc2x0 C. Other sulfonate esters can also be used (e.g., tosylate, nosylate or mesylate) as well as halides such as epibromohydrin or epichlorohydrin. The nosylate is preferred. Specific reaction conditions are described in Examples 1 and 2.\nThe hydrolysis of the indazolyloxy propanolamine ester represented by Structural Formula (XIII) to form the carboxylic acid represented by Structural Formula (XVI) or the carboxylate salt thereof can be carried out by any suitable means, including by procedures disclosed in Larock, R. C. Comprehensive Organic Transformations; VCH: New York, 1989, pp. 981-985; March, xe2x80x9cAdvanced Organic Chemistryxe2x80x9d, third(3rd) edition, John Wiley Sons (1985), pages 375-76 and references cited therein, the entire relevant teachings of which are incorporated herein by reference. This reaction is referred to herein as the xe2x80x9chydrolysis reactionxe2x80x9d.\nPreferably, the hydrolysis reaction is carried out in alcoholic base such as lithium hydroxide, sodium hydroxide or potassium hydroxide in methanol or ethanol. Methanolic sodium hydroxide is preferred. The hydrolysis can be carried out with one equivalent of base relative to indazolyloxy propanolamine ester. Alternatively, an excess of base, for example, an excess up to about ten to about twenty-fold can be used to form the carboxylate salt. Preferably, the excess of base is between about zero and about 30 percent. Generally, the reaction temperature varies between room temperature and the reflux temperature of the solvent, and is typically 40-70xc2x0 C. The carboxylic acid product (or carboxylate salt thereof) is isolated by conventional means, for example, by removal of the solvent in vacuo.\nThe indazolyloxy propanolamine carboxylic acid or carboxylate salt thereof can be purified by basic (or acidic) extraction, e.g., by dissolving in aqueous base (e.g., aqueous sodium hydroxide, potassium hydroxide or lithium hydroxide), and/or by precipitation at its isoelectric point. Preferably, the aqueous solution is extracted with one or more organic solvents which are not miscible with water to remove impurities. The pH of the aqueous solution can then adjusted to its isoelectric point with aqueous acid (e.g., aqueous HCl, H2SO4, acetic acid or a sulfonic acid), thereby causing the indazolyloxy propanolamine carboxylic acid to precipitate.\nThe indazolyloxy propanolamine carboxylic acids or carboxylate salts thereof can also be purified by using either cationic ion-exchange resins, such as AMB 15 (H form), 50WX2-400 (H form), or IR50S (H form) and eluting with an alcoholic solution of an alkali metal acetate, such as 6% (w/v) sodium acetate in methanol; or by using anionic ion-exchange resins, such as IRA900 (chloride form) or A21 (acetate form) and eluting with a basic solution of water and either alcohol or acetonitrile.\nThe individual optically active isomers of the compounds prepared by the present invention may be prepared from their respective optically active precursors by the procedures described above, or by resolving the racemic mixtures. This resolution can be carried out by derivatization with a chiral reagent followed by chromatography or by repeated crystallization. Removal of the chiral auxiliary by standard methods affords substantially optically pure isomers of the compounds of the present invention or their precursors. Further details regarding resolutions can be obtained in Jacques, et al., Enantiomers, Racemates, and Resolutions, John Wiley Sons, 1981.\nThe compounds employed as initial starting materials in the synthesis of the compounds of this invention are well known and, to the extent not commercially available, are readily synthesized by standard procedures commonly employed by those of ordinary skill in the art.\nThe invention is illustrated by the following examples, which are not meant to be limiting in any way."}
{"text": "MicroElectroMechanical System (MEMS) assemblies include microphones and speakers to mention two examples. These MEMS devices may be used in diverse applications such as within hearing aids and cellular phones.\nIn the case of a MEMS microphone, acoustic energy typically enters through a sound port in the assembly, vibrates a diaphragm and this action creates a corresponding change in electrical potential (voltage) between the diaphragm and a back plate disposed near the diaphragm. This voltage represents the acoustic energy that has been received. Typically, the voltage signal is then transmitted to an electric circuit (e.g., an integrated circuit such as an application specific integrated circuit (ASIC)). Further processing of the signal may be performed on the electrical circuit. For instance, amplification or filtering functions may be performed on the voltage signal by the integrated circuit.\nAs mentioned, sound typically enters the assembly through an opening or port. When a port is used, this opening also allows other unwanted or undesirable items to enter the port. For example, various types of contaminants (e.g., solder, flux, dust, and spit, to mention a few possible examples) may enter through the port. Once these items enter the assembly, they may damage the internal components of the assembly such as the MEMS device and the integrated circuit.\nPrevious systems have sometimes deployed particulate filters that prevent some types of debris from entering an assembly. Unfortunately, these filters tend to adversely impact the operation of the microphone. For instance, the performance of the microphone sometimes becomes significantly degraded when using these previous approaches. Microphone customers often elect to not use such microphones in their applications because of the degraded performance.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity. It will further be appreciated that certain actions and/or steps may be described or depicted in a particular order of occurrence while those skilled in the art will understand that such specificity with respect to sequence is not necessarily required. It will also be understood that the terms and expressions used herein have the ordinary meaning as is accorded to such terms and expressions with respect to their corresponding respective areas of inquiry and study except where specific meanings have otherwise been set forth herein."}
{"text": "A headphone having a noise-cancelling function (noise-cancelling headphone) includes a microphone and a noise-cancelling circuit (hereinafter referred to as “NC circuit”). The microphone collects external sounds (hereinafter referred to as “noises”) around the headphone. The NC circuit generates cancelling signals corresponding to the noises collected by the microphone. The headphone combines sound waves corresponding to the cancelling signals generated by the NC circuit and sound waves corresponding to the reproduced signals from a sound source such as a music player connected to the headphone, and then outputs the combined sound waves from a driver unit. That is, the headphone outputs musical sounds (hereinafter referred to as “reproduced sounds”) derived from the reproduced signals from the sound source with cancelling (muting) of the noises.\nExamples of the scheme for generating cancelling signals include a feedback scheme (hereinafter referred to as FB) and a feedforward scheme (hereinafter referred to as FF).\nA built-in microphone of the FB noise-cancelling headphone is disposed in the interior of a housing unit (ear piece) of the headphone and near the ear of the user. The NC circuit analyses in real-time the signals of the noises collected by the microphone and generates cancelling signals minimizing the noises at the position of the eardrum of the user. The FB noise-cancelling headphone collects the noises at the position near the ear of the user. Thus, the noise-cancelling effect of the FB noise-cancelling headphone is higher than that of the FF noise-cancelling headphone. In addition, the response of the FB noise-cancelling headphone to the variation of noise components is good.\nHowever, when the built-in microphone of the FB noise-cancelling headphone collects the reproduced sounds in addition to the noises, and then the NC circuit generates cancelling signals, the sound quality of the reproduced sounds output from the headphone are degraded. In addition, the ear piece of the FB noise-cancelling headphone should be sealed in the state which the headphone is worn on the head of the user, to enhance the noise-cancelling effect. When the ear piece is sealed, the reproduced sounds output from the headphone may be muffled. Thus, the sound quality of the reproduced sounds are degraded. To address the problem, a FB noise-cancelling headphone generally includes a filter to correct the sound quality of the reproduced sounds.\nOn the other hand, a built-in microphone of the FF noise-cancelling headphone is disposed at the exterior of the housing unit of the headphone. The NC circuit analyses the signals of the noises collected by the microphone and predicts the variation in the noises that will reach the eardrum of the user wearing the headphone. The NC circuit generates cancelling signals based on the result of the prediction. The FF noise-cancelling headphone does not need the placement of the microphone on a limited space near the ear of the user. In addition, the built-in microphone of the FF noise-cancelling headphone is disposed at a position remote from the driver unit. Thus, the FF noise-cancelling headphone is less likely to collect the reproduced sound output from the driver unit and to generate cancelling signals from the reproduced sounds than the FB noise-cancelling headphone. That is, the sound quality of the reproduced sound output from the FF noise-cancelling headphone are less susceptible to the cancelling signals generated from the reproduced sounds than the sound quality of the reproduced sound output from the FB noise-cancelling headphone.\nUnfortunately, the FF noise-cancelling headphone often indicates directionality of the noise-cancelling effect depending on the position of the built-in microphone. Unlike the built-in microphone of the FB noise-cancelling headphone, the built-in microphone of the FF noise-cancelling headphone is disposed at the exterior of the housing unit. Thus, the FF noise-cancelling headphone generates noises caused by wind pressure of the blowing of the wind, for example. In other words, the FF noise-cancelling headphone is susceptible to the influence of the wind around the headphone. As a result, the FF noise-cancelling headphone can cause a feeling of strangeness or discomfort to the user, when the noise-cancelling function is activated.\nSchemes have been proposed to dispose a microphone in the interior of the housing unit of the FF noise-cancelling headphone to avoid the influence of the wind described above (for example, refer to Japanese Unexamined Patent Application No. 2010-109799).\nThe FF noise-cancelling headphone disclosed in Japanese Unexamined Patent Application No. 2010-109799 includes a driver unit, a baffle plate, and a microphone. The microphone is disposed behind the baffle plate to which the driver unit is attached (that is, in the interior of a rear air chamber). The baffle plate includes a flange portion and a groove. The flange portion is disposed on the forward portion and the rear portion in the thickness direction of the baffle plate. The groove is formed between the flange portion and along the outer circumference of the baffle plate. The flange portion that is disposed on the rear portion of the baffle plate has a sound collecting hole penetrating the flange portion in thickness direction. The sound collecting hole is in communication with the groove. The external noises are collected by the microphone through the groove and the sound collecting hole."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrophotographic photoconductor that is used in electrophotographic printers, copiers, fax machines and the like (hereafter also referred to simply as “photoconductor”), to a method for producing the electrophotographic photoconductor, and to an electrophotographic apparatus. More particularly, the present invention relates to an electrophotographic photoconductor that, by containing a polymer of specific structure, exhibits excellent contamination resistance, electrical characteristic stability, and ozone resistance, to a method for producing the electrophotographic photoconductor, and to an electrophotographic apparatus.\n2. Background of the Related Art\nThe electrophotographic photoconductor has a basic structure wherein a photoconductive layer having a photoconductive function is disposed on a conductive support. Ongoing research and development are actively carried out on organic electrophotographic photoconductors that utilize organic compounds as functional components for charge generation and transport, given the advantages that organic electrophotographic photoconductors afford, such as material diversity, high productivity, safety and the like. The use of organic electrophotographic photoconductors in copiers, printers and the like is thus becoming more widespread.\nGenerally, a photoconductor must fulfill the function of holding surface charge in the dark, generating charge when receiving light, and transporting the charge thus generated. Such photoconductors encompass so-called single layer-type photoconductors that are provided with a single photoconductive layer that combines the above functions, and so-called multilayer-type (of function-separated type) photoconductors provided with a photoconductive layer that is a stack of layers functionally separated into a charge generation layer that fulfils mainly the function of charge generation upon reception light, and a charge transport layer that fulfils the function of transporting the charge that is generated in the charge generation layer upon reception of light.\nThe above photoconductive layers are generally formed by coating a conductive support with a coating solution in which a charge generation material, a charge transport material and a binder resin are dissolved or dispersed in an organic solvent. In many instances a polycarbonate is used as the binder resin in organic electrophotographic photoconductors, in particular in the outermost surface layer, since polycarbonates are resistant to friction with paper and with blades for toner removal, and boast excellent flexibility and good exposure transparency. Among the foregoing, bisphenol Z polycarbonates are widely used as the binder resin. Technologies in which such polycarbonates are used as a binder resin include, for instance, Japanese Patent Application Publication No. S61-62040 (Patent literature 1).\nElectrophotographic printing devices are required to possess ever higher durability and sensitivity, and faster responses, to cope with, for instance, increases in the number of prints to be printed in a networked office, and with the rapid development of lightweight electrophotographic printing machines. These devices, moreover, are held to strict requirements in terms of being little affected by gases, such as ozone and NOx that are generated in the device, and exhibiting little fluctuation in image characteristics arising from variations in the usage environment (room temperature and humidity).\nRecent developments in color printers, and the growing prevalence of the latter, have been accompanied with a need for higher printing speeds, smaller equipment and fewer constituent members, as well as the need for accommodating various usage environments. Under such circumstances, there is a pressing demand for photoconductors that exhibit little variation in image characteristics and electrical characteristics caused by repeated use and/or derived from fluctuations in the usage environment (room temperature and environment). Conventional technologies have thus far failed to meet these requirements simultaneously to a sufficient degree.\nOzone is a widely known example of a gas that is generated in equipment. Ozone is generated by a charger or roller charger that triggers corona discharge. The photoconductor becomes thus exposed to residual ozone or dwelling ozone within the equipment. It is found that the organic substances that make up the photoconductor become oxidized thereby, and, as a result, the original structure of the photoconductor breaks down, and the photoconductor characteristics are significantly impaired. Moreover, it is found that ozone oxidizes the nitrogen in air into NOx, which in turn alters the organic substances that make up the photoconductor.\nIt is deemed that not only does characteristic deterioration elicited by such gases extend to the outermost layer of the photoconductor, but also adverse effects arise when the gas flows into the interior of the photoconductive layer. It is found that the outermost layer itself of the photoconductor is scraped off, though slightly, on account of friction with the above-described various members. When a harmful gas flows into the interior of the photoconductive layer, the organic substances in the photoconductive layer may undergo structural breakdown. Suppressing the inflow of such harmful gas is thus an issue to be addressed. In tandem-type color electrophotographic apparatuses that rely on a plurality of photoconductors, in particular, variation in color tone occurs as a result of differences in the degree of influence of the gas, depending on, for instance, the position at which drums are disposed in the device. Such variations are deemed to constitute an impediment to forming adequate images. Therefore, it is found that characteristic deterioration caused by gas is a particularly important issue in tandem-type color electrophotographic apparatuses.\nThe photoconductor surface may also be contaminated by ozone, nitrogen oxides and the like that are generated during charging of the photoconductor. Problems that arise in such a case include, for instance, image smearing by those contaminants themselves, as well as lowered surface lubricity caused by the adhered substances, greater likelihood of adhesion of paper dust and toner, and likelier occurrence of blade squealing, curling, surface scratches and the like.\nIt is also found that human sebum and the like becomes adhered to the photoconductor surface during repair of the electrophotographic apparatus and during the operation of replacing photoconductor units. Thus far, however, the durability of photoconductors against such contaminant adhesion has been not necessarily sufficient, and surface cracks, as well as image defects such as white spots and black spots occur in some instances when human nose fat or scalp sebum is left adhered to the surface of the photoconductor over long periods of time.\nVarious methods for improving the outermost surface layer of photoconductors have been proposed in order to solve the above problems. Specifically, various polycarbonate resin structures have been proposed in order to enhance the durability of photoconductor surfaces. For instance, Japanese Patent Application Publication No. 2004-354759 (Patent literature 2) and Japanese Patent Application Publication No. H04-179961 (Patent literature 3) propose polycarbonate resins that comprise a specific structure, but not enough consideration is given to compatibility with various charge transport agents and various additives, or to resin solubility. For instance, Japanese Patent Application Publication No. 2004-85644 (Patent literature 4) proposes a polycarbonate resin that comprises a specific structure; however, a resin having a highly bulky structure includes large spaces between polymers, and thus substances released during charging, as well as contact members and foreign matter, permeate readily into the photoconductive layer, and it is accordingly difficult to achieve sufficient durability. Japanese Patent Application Publication No. H03-273256 (Patent literature 5) proposes a polycarbonate having a special structure, in order to enhance printing durability and coatability, but does not sufficiently disclose additives or charge transport materials that are combined with the polycarbonate. Patent literature 5 is problematic in that maintaining stable electrical characteristics over long periods of time is difficult.\nJapanese Patent Application Publication No. 2010-276699 (Patent literature 6) proposes the feature of adding a highly branched polymer and a polymerizable charge transport agent to a surface protective layer, to enhance thereby abrasion properties and transfer properties, but coating solution stability is still an issue. Japanese Patent Application Publication No. 2003-255580 (Patent literature 7) proposes the feature of incorporating, into the surface layer of the photoconductor, a binder resin and a linear vinyl polymer having long-chain alkyl groups at side chains. However, it is found that upon polymerization of a vinyl polymer in a solution in the presence of another binder resin, it is difficult to control the molecular weight and the resin skeleton, due to presence of that other resin. Regarding improvements in wear resistance by a surface protective layer, Japanese Patent Application Publication No. 2011-64734 (Patent literature 8) proposes a technology that involves configuring a surface protective layer that contains a cured product having a three-dimensional crosslinked structure and formed out of a predetermined radically polymerizable compound, a trifunctional or higher functional radically polymerizable monomer, and a radically polymerizable compound having a charge transporting structure, but this configuration is problematic in terms of productivity, since the photoconductive layer has a multilayer structure. Improvements derived from the charge transport layer are an important issue herein.\nTransfer current tends to increase in color printers on account of toner color overlap and/or the use of transfer belts. When printing on paper of various sizes, a difference in transfer fatigue arises between portions with paper and portions without paper. This in turn exacerbates differences in image density, which is problematic. In case of frequent printing on small-sized paper, bare photoconductor portions over which the paper does not pass (paper non-passage sections) are continuously and directly affected by transfer, and exhibit greater transfer fatigue than bared photoconductor portions over which paper does pass (paper passage sections). As a result, when printing is subsequently performed on large-size paper, the above discrepancy in transfer fatigue between paper passage sections and paper non-passage sections gives rise to a potential difference in the developed area, which translates into observable differences in density. This trend becomes yet more pronounced as transfer current increases. Under such circumstances, the demand has intensified for photoconductors that exhibit little fluctuation in image characteristics and electric characteristics as a result of repeated use, or on account of fluctuations in the usage environment (room temperature and environment), and that exhibit excellent transfer resiliency, particularly in color printer, as compared to monochrome printers. Conventional technologies have thus far failed to meet these requirements simultaneously to a sufficient degree.\nVarious additives, such as hindered phenol compounds, phosphorus compounds, sulfur compounds, amine compounds, hindered amine compounds and the like have been proposed to enhance gas resistance. The current situation, however, is that these technologies fail to provide sufficient gas resistance, or, even if satisfactory characteristics are exhibited in terms of gas resistance, no satisfactory results are achieved, through a combination of resins and charge transport materials, regarding electrical characteristics, for instance, responsiveness, image memory and potential stability in endurance printing. The applicants had proposed diester compounds in WO 2011/108064 (Patent literature 9) and Japanese Patent Application Publication No. 2007-279446 (Patent literature 10), but have since made further progress in the study of combinations of more appropriate binder resins and high-mobility charge transport materials.\nVarious conventional technologies have been proposed pertaining to improvement of the surface layer of photoconductors. However, these technologies as disclosed in the citations above were not all sufficient as regards electrical characteristics such as light response, and also contamination resistance towards sebum, photoconductor productivity and the like.\nTherefore, it is an object of the present invention to provide an electrophotographic photoconductor that has excellent contamination resistance and stable electrical characteristics and so forth upon repeated use, and superior transfer resistance and gas resistance, and to provide a method for producing the electrophotographic photoconductor, and an electrophotographic apparatus."}
{"text": "1. Field of the Invention\nThe field of invention relates to pillow construction, and more particularly pertains to a new and improved pillow template apparatus wherein the same is arranged for the severing of a fabric web for pillow construction while simultaneously orienting a pattern relative to the fabric web.\n2. Description of the Prior Art\nThe cutting and removal of a web portion for construction of a pillow results in a non-uniform typical construction of the resultant pillow should the individual severing such a web not accommodate and provide for a convex side wall construction, whereupon stuffing the pillow permits a final resultant product of linear size. While the prior art, such as exemplified in U.S. Pat. Nos. 3,440,980; 4,930,382; 3,774,558; and 5,058,285 provides for various template structure relative to sewing procedures, the instant invention overcomes deficiencies of the prior art by providing for a template structure arranged for the specific severing of a web for the construction of a symmetrically oriented pillow and in this respect, the present invention substantially fulfills this need."}
{"text": "Optical attenuation of interferers in a radio frequency signal has been accomplished by modulating a beam of coherent light with the signal, removing the unwanted interference spectral components from the modulated beam, and then down-converting the beam to provide a radio signal which is free of these interference components.\nSuch a technique is shown in U.S. Pat. Nos. 4,699,466 and 4,522,466, assigned to the present assignee. The teachings of these patents are meant to be incorporated herein by way of reference.\nThe aformentioned technique filters the modulated radiation through an optical Fourier transform lens, a spatial filter and an inverse Fourier transform lens.\nIn particular, U.S. Pat. No. 4,522,466 teaches that the filter attenuation can be maximized by utilizing recursive methods. The coherent light is passed through the optical filtering system a number of times to provide extremely high notching attenuation on the order of 40 dB, or greater.\nIt would be desirable, however, to increase the notch depth to 80 dB, or more, or possibly notch the signal to the noise floor to remove all EMI.\nThe invention has as one of its objectives to provide an improved system and method of optically processing radio frequency signals to remove unwanted EMI, wherein a greater notch depth can be achieved than has been heretofore accomplished.\nAnother object of this invention is to provide a technique and system which reduces the number of filtering elements in a cascaded, multiple channel optical array."}
{"text": "1. Field of the Invention\nThe present invention relates to a water dissolvable or water dispersible produce label for individually labeling various fruits and vegetables and more specifically, to a dissolvable produce label that contains an organic cleansing produce wash to help clean the individually labeled fruit or vegetable by aiding in the removal of wax, dirt, bacteria, pesticides and fungicides. This invention further relates to a dissolvable produce-cleansing produce label that is water resistant in transit.\n2. Background\nCurrent produce labels require removal by the end user. The end user must first peel the produce label off of each fruit individually and dispose of the label(s) before washing and/or eating. These labels are seen as a source of frustration by the end user. They are, for the most part, difficult to peel off and separate from your finger during disposal and add waste to an ever-growing trash problem. Also, current produce labels do nothing in the way of helping the user to clean their fruit. In fact, current fruit labels leave behind an adhesive residue that is very difficult to wash off with water alone.\nThe end user is warned against using detergent, laundry bleach or soap to clean their produce, as fruits and vegetables are porous and can absorb the soap or bleach. Therefore, produce washes have been formulated to help clean fruits and vegetables better than washing with water alone. There are currently dozens of companies manufacturing and selling “produce wash” as an expensive aftermarket cleansing solution to help consumers clean their fruits and vegetables. Manufacturers currently producing Produce Wash include: Veggie Wash, Eat Cleaner Fruit & Vegetable Wash, Fit Fruit & Vegetable Wash, Ultimate Fruit & Vegetable Wash, Mom's Veggiewash, Earth Friendly Fruit & Vegetable Wash, Enviro-one Fruit Wash, etc. . . . .\nProduce labeling provides for fast and accurate checkout. In addition to offering basic product information, the PLU codes printed on produce labels also indicate whether a fruit is conventional, organic or genetically modified. Some labels now include a HarvestMark Number that provides additional information like where it was grown and whether or not the product is subject to a recall. With the HarvestMark number, the end user can also view photos of the farm, or the story behind the farmer, give feedback about the product they purchased, and discover nutrition information, and even recipes. Produce labels also provide branding space to build consumer trust. With branding in mind, fruit growers, distributors and supermarket chains are continually looking for opportunities to responsibly show their commitment to providing their customers with a clean product and a more enjoyable user experience.\nIt is desirable to have a fruit label that can be washed away; that leaves no adhesive residue behind; and that helps the end user clean their fruit. It is also desirable to offer a dissolvable fruit label that may be removed by “peeling it off” should the user choose to do so."}
{"text": "Conventional display devices often have the defect of display unevenness (e.g., sometimes referred to as “Mura” or “Trace Mura”) in a part of display regions, Mura is a poor display severely affecting the quality of the picture, the severity and duration of Mura are related to many factors, one is process parameter of the panel, such as liquid crystal, orientation layer, and the like, and the other is the setting of driving conditions. Experimental results show that white state voltage of the liquid crystal is reduced to effectively reduce the severity of the Mura and shorten its duration, wherein the white state voltage refers to the driving voltage of the display device corresponding to the maximal brightness.\nThe white state voltage of the liquid crystal is reduced to improve the Mura, but certain negative effects shall also be brought. As shown in FIG. 1, the transmittance shall be reduced due to the reduction of the white state voltage, namely, the maximal brightness is reduced, and meanwhile, the reduction of the maximal brightness means contrast reduction. With the reduction of the white state voltage, the variation of the driving voltage of each average grey scale is decreased, so that more accurate driving voltage of the display device is needed, which makes a driver integrated circuit (Driver IC) face more severe working conditions.\nTo sum up, in the prior art, the white state voltage of the liquid crystal is reduced to improve the Mura, but the maximal brightness shall be reduced, and meanwhile, the driver integrated circuit will face more severe working conditions."}
{"text": "The use of a CATV system to provide internet, voice over internet protocol (“VOIP”) telephone, television, security, and music services is well known in the art. In providing these services, a downstream bandwidth (i.e., radio frequency (“RF”) signals, digital signals, and/or optical signals) is passed from a supplier of the services to a user, and an upstream bandwidth (i.e., radio frequency (“RF”) signals, digital signals, and/or optical signals) is passed from the user to the supplier. For much of the distance between the supplier and the user, the downstream bandwidth and the upstream bandwidth make up a total bandwidth that is passed via a signal transmission line, such as a coaxial cable. The downstream bandwidth is, for example, signals that are relatively higher in frequency within the total bandwidth of the CATV system while the upstream bandwidth is, for example, signals that are relatively lower in frequency.\nTraditionally, the CATV system includes a head end facility, where the downstream bandwidth is initiated into a main CATV distribution system, which typically includes a plurality of trunk lines, each serving at least one local distribution network. In turn, the downstream bandwidth is passed to a relatively small number (e.g., approximately 100 to 500) of users associated with a particular local distribution network. Devices, such as high-pass filters, are positioned at various points within the CATV system to ensure the orderly flow of downstream bandwidth from the head end facility, through the trunk lines, through the local distribution networks, and ultimately to the users.\nAt various locations between the head end facility and the user, there are amplifiers and slope adjustment devices for the purpose of maintaining the quality of the downstream bandwidth. This statement introduces three terms (i.e., quality, amplifiers, and slope adjustment devices) that are important to the remaining discussion. These will be discussed broadly below.\nThe quality of the downstream bandwidth is often a measure of: (i) a signal level of a particular channel within the downstream bandwidth, the signal level referred to merely as “level;” and (ii) a general consistency of levels across all of the channels in the downstream bandwidth, the general consistency referred to as “slope.” These objective measurements are often used by technicians to evaluate CATV system performance during operation and to troubleshoot customer complaints.\nThe level of each channel should fall within a specific range that has been determined to provide satisfactory video, sound and information transfer rates for users. The specific requirements for each channel are not of importance to the present discussion, but it is helpful to understand that are specific targets for the level of each channel. Note that this is a simplistic definition to explain “level,” and note that this definition does not include other variances such as between analog and digital.\nSlope is measurement used to assess the amount of loss experienced due in large part to cable length. While all channels experience some loss, channels transmitted using higher frequencies within the downstream bandwidth experience more loss than those transmitted using lower frequencies. Accordingly, when the levels for all of the channels within the downstream bandwidth are graphed such that they are arranged in order according to the frequency of the channel, there may be a significant visual downward slope in the graph from the lowest frequency channel to highest frequency channel. This downward slope becomes more prominent as the length of signal cable increases. Note that this is a simplistic definition to explain the consistency of levels across all of the channels and the “slope” that is created by losses occurring in the signal cables. Also note that this definition does not include other variances such as between analog and digital.\nThe presence of slope is not removed through the use of typical drop-style amplifiers. The drop-style amplifiers merely amplify the entire downstream bandwidth. In other words, these drop-style amplifiers raise the level of each channel equally. In turn, if there is a large amount of slope present, such as when a user's premise includes long distances of signal cable, the drop-style amplifier may cause some channels to exceed their level specification while other channels may remain below their specification.\nIt is known to add a fixed or manually adjustable slope compensator/low frequency attenuator when there is a long run of signal cable. However, these devices require expensive testing equipment to determine whether and/or how much slope compensation should be supplied to a particular premise. Further, due to the cost of installation and a general misunderstanding regarding how to install such devices, there are relatively few in existence, compared to the number of such devices needed.\nFor at least the forgoing reasons, a need is apparent for a device, which can increase the overall quality of the downstream bandwidth that includes providing an appropriate amount of amplification and an appropriate amount of slope compensation increasing the signal strength and increasing the signal-to-noise ratio."}
{"text": "The popularity of multiplayer online games on the Internet has exploded in recent years. Unlike conventional single-player electronic games, multiplayer online games enable a player to participate in games with other players, either as an individual or as a team member. The multiplayer aspect of these games adds to the players' gaming experience by offering new gaming considerations, including strategy, teamwork, and rivalries.\nIn order to enable playing of a multiplayer online game, host services must be provided by either an individual or a third-party host. Host services are typically performed on a host computer or server running a copy of the game, and generally concern controlling overall aspects of the game, such as keeping track of the score and updating the game environment. In addition, the host determines who the game participants are, and whether new players may be added.\nExamples of multiplayer online games that are hosted by third-party providers can be found at many gaming web sites, including Microsoft Corporation's Internet Game Zone, which provides a myriad of single player and multiplayer online games that are played by over a half-million users daily. In order to play most of the games on a gaming site, it is necessary for users to download (or procure separately) a copy of the game the user wishes to play. Typically, some of these games may be downloaded for free, while others must be purchased, either online, or at a retail outlet. In addition, many gaming sites offer “premium” games that can only be accessed if the user has paid a membership fee to play such games, or has agreed to pay a usage fee (e.g., an hourly fee, daily fee, etc.).\nIn addition to Internet gaming site hosts, individuals may also provide hosting services by using their own computer for such purposes. Such a player is referred to as a “host player.” A significant advantage of this configuration is that it generally provides increased performance over games hosted by third-party sites, since the host services for these sites are usually provided by a limited number of servers, each of which is used to simultaneously host a large number of games. In addition, by using one of the game players as the host for the game, the players do not have to pay an access charge to play the game.\nIn order to facilitate online game playing, it is necessary to enable players to join games already in progress, or enable players to initiate a new game in a manner that allows other players to join in. Oftentimes, players prefer to play games against (or with) players they already know. In general, many of the third-party gaming sites provide host services that make it relatively easy to join games, schedule future games, establish teams, etc. However, there is presently no way to easily initiate game play between players for games that are hosted by individual players.\nIn order to host a game, the host needs to be able to identify and communicate with all of the game participants. A common scheme for performing this task involves the use of Internet Protocol (IP) addressing. In order to communicate over a communications network, such as the Internet, each device (e.g., a player's computer) connected to the network must be assigned to a unique network address. Under the Internet's present TCP/IP (Transmission Control Protocol/Internet Protocol) scheme (which is also used for other networks), each device connected to the Internet is identified by a 32-bit IP address comprising four 8-bit segments separated by dots, e.g., 207.21.32.212. Each IP address comprises a unique network ID, and a unique host ID, the latter of which identifies a workstation, server, router, or other TCP/IP device on the network. Accordingly, each of the individual computers used by the game players is assigned to its own IP address. As a result, a game host can enable new players to join a game by telling the other players the IP address of the host's machine.\nInitiating a game in this matter is burdensome and prone to error. First, all of the players must be apprised of the host's IP address. In general, the host must inform each player of the host's IP address, unless the host uses a computer that is assigned to a static (i.e., permanent) IP address. For example, most game players connect to the Internet via a dial-up online Internet Service Provider (ISP), such as AMERICA ONLINE™, MSN™, COMPUSERVE™, EARTHLINK™, MINDSCAPE™, etc. Each of these ISPs has a fixed number of connections it can use simultaneously, with each connection being assigned to a unique, static IP address. In general, the number of connections (and thus IP addresses) is only a fraction of the number of customers to whom the ISP provides service. Accordingly, the ISPs use dynamic IP addressing schemes, whereby users are dynamically assigned to a connection (and the connection's corresponding IP address) upon logging onto the ISP, and the IP address of the connection is released for use by a subsequent user when the initial user logs off of the ISP. As a result, there is no way to assign a static IP address to users of ISPs of this type. Accordingly, the only way the players can be apprised of the host's IP address is for the host to let them know what the IP address is after it has been assigned, which is typically done through e-mail or via a phone call.\nAn improvement over the foregoing manual method is provided by online lookup services. In this instance, each member of the online lookup service is assigned a unique alias (i.e., user name) upon signing up for the service, wherein that user is always identified by the alias whenever the user logs onto the lookup service. In addition, upon logging on, the IP address for the user is matched to the user's alias so that the user's IP address can be accessed by application programs running on other computers through use of the alias. As a result, players can then search the online lookup service for the aliases of a host or hosts the players know are or will be hosting a game, wherein the IP addresses of those hosts is passed to the players (generally behind the scenes to an application program running on the player's computer). At this point, the players may still need to manually enter the IP address of the host in order to join the game.\nAlthough the online lookup service scheme is an improvement over the manual method, it is still deficient in many respects. Notably, this scheme doesn't provide a convenient way to invite players to join games, wherein the invited players can decide whether they want to join the game. Additionally, the host must access multiple application programs and perform many steps in order to initiate a multiplayer game. Accordingly, it is desired to provide an improved scheme for enabling individuals to host multiplayer computer games and other players to join such games that overcomes the many limitations of the foregoing manual and online lookup service schemes."}
{"text": "Floating offshore structures such as e.g. oil platforms, offshore wind turbines, and offshore meteorological towers are known.\nSeveral configurations have been proposed for the floating or buoyancy structures: many of these employ floater elements in the form of substantially hollow floater tanks that in use are arranged substantially below the mean sea level and provide a buoyancy force to support the structure, e.g. a wind turbine. Ballast and/or mooring lines anchored to the seabed are provided for achieving stability.\nIn some of these floating structures, the buoyancy structure is designed to provide an excess buoyancy force and is maintained floating under the mean sea level by taut mooring lines tensioned by the excess buoyancy force.\nFor example, concepts have been developed such as the “Taught Leg Buoy” (TLB) floating wind turbine, with a slender cylindrical buoy and two sets of tensioned mooring lines, inclined relative to the seabed and connected to gravity anchors and to the buoy; or such as the “Tension Leg Platform” (TLP) floating wind turbine, in which the tensioned mooring lines are substantially vertical and are connected between gravity anchors on the seabed and arms or braces extending radially outwards with respect to the vertical axis of the wind turbine. The TLP arms may be part of the buoyancy structure, for example in the form of hollow spokes that extend radially outward from a hollow central hub, or may be arranged above the sea level, in which case the buoy may be a slender cylindrical tank like in the TLB concept.\nTLP structures are also used for offshore oil platforms. Other configurations for offshore structures are not based on providing an excess buoyancy force and tensioned mooring lines for their stability, but instead on having a ballast at a relatively point of the structure. The resulting centre of gravity may thus be located at a relatively low point of the structure as well. As such, the resulting structure may be inherently stable and does not depend on the tensioned mooring lines for stability.\nThe buoyancy structures of offshore structures are subject to several loads, such as for example the weight of the wind turbine itself, impacts, forces exerted by waves, currents and tides, and, in case of a wind turbine, also aerodynamic forces associated to the wind, rotor rotation, etc.\nAs mentioned before, the offshore structures that rely on excess buoyancy for stability have tensioned mooring lines or cables. In order for these structures to be stable, the mooring lines must always be under tension, otherwise the offshore structure could become unstable and could fall over. Since the loads on the offshore structures, and particularly on a floating wind turbine may vary considerably both in magnitude and direction, a high amount of excess buoyancy and high tension in the cables may be required.\nHaving a high tension in the mooring lines or cables means that the cross-sectional area of the cables needs to be increased as well. The resulting cables may thus be heavy and expensive. Additionally, having such a high tension in the cables may result in the rest of the structure to be dimensioned accordingly. Also the other parts of the structure may become heavier and more expensive.\nAn alternative solution would be to have a high number of lines for providing stability and distributed in such a way that even if some of them could in certain circumstances enter in compression, the other lines would still be able to maintain the stability of the offshore structure. However, also this solution can become very expensive.\nThere thus exists a need for a floating offshore structure that reduces or resolves at least some of the afore-mentioned problems."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a program register using a nonvolatile memory device and a programmable logic circuit using the same, and more specifically, to a technology for storing data or performing an operation on the data without additional memory devices, thereby reducing the area of the circuit.\n2. Description of the Prior Art\nGenerally, a ferroelectric random access memory (hereinafter, referred to as ‘FRAM’) has attracted considerable attention as next generation memory device because it has a data processing speed as fast as a Dynamic Random Access Memory DRAM and conserves data even after the power is turned off.\nThe FRAM having structures similar to the DRAM includes the capacitors made of a ferroelectric substance, so that it utilizes the characteristic of a high residual polarization of the ferroelectric substance in which data is not deleted even after an electric field is eliminated.\nThe technical contents on the above FRAM are disclosed in the Korean Patent Application No. 1999-49972 by the same inventor of the present invention. Therefore, the basic structure and the operation on the FRAM are not described herein.\nA conventional programmable logic operation circuit for changing logic levels of input signals stores address information in storage means. However, since a SRAM (Static Random Access Memory) is used as the conventional programmable logic operation circuit, various information stored in latches is leaked in a power-off mode. Even when power is supplied to the system again, various data for operations of circuits are required to be reset."}
{"text": "This invention relates to radiation curable compositions. More particularly, this invention relates to radiation curable compositions comprising an alkenyl ether functional polyisobutylene, a cationic photoinitiator and a miscible reactive diluent. This invention can further comprise of optional ingredients such as free radical photoinitiator, photosensitizers, non-reactive diluents, thermo-oxidative stabilizers and shelf life stabilizers.\nPolyisobutylenes containing functional groups which are radiation curable have been disclosed in the art. For example, T. P. Liao and J. P. Kennedy in Polymer Bulletin, V. 6, pp. 135-141 (1981) disclose acryl and methacryl telechelic polyisobutylenes having the formula CH2xe2x95x90C(R)xe2x80x94COO-PIB-OOCxe2x80x94C(R)xe2x95x90CH2 where R is xe2x80x94H or CH3. These materials were prepared by reacting alpha, omega di-hydroxypolyisobutylene, HOCH2-PIB-CH2OH, and excess acryloyl or methacryloyl chloride. These prepolymers are disclosed as being useful in the synthesis of a variety of new composites containing a soft polyisobutylene segment.\nJ. P. Kennedy and B. Ivan in Polymer Material Science and Engineering, V. 58, p.866 (1988) disclose allyl telechelic linear and star-branched polyisobutylenes prepared by a convenient rapid one pot polymerization functionalization process. The polymerization step involved living polymerization of isobutylene by recently discovered mono- or multifunctional initiating systems (combinations of tert.-ester and ether/Lewis acids) followed by electrophilic functionalizations by allyl trimethylsilane in the presence of TiCl4. Characterization indicated quantitative end allylations. Subsequent quantitative derivations of the allyl termini yielded mono-, di-, and tri-functional hydroxyl- and epoxy-telechelic polyisobutylenes which could be cured to rubbery networks.\nJ. P. Kennedy and B. Ivan in the Journal of Polymer Science, Part A, Polymer Chemistry, V. 28, p. 89 (1990) disclose mono-, di-ended linear, and three-arm star allyl telechelic polyisobutylenes which are prepared by a rapid economical one-pot polymerization-functionalization process. The process involved the living polymerization of isobutylene by mono-, di-, or tri-functional initiating systems, specifically by aliphatic and aromatic tert-ester and -ether/TiCl4 combinations, followed by electrophilic functionalization of the living sites with allyl-trimethylsilane. Quantitative derivations of the ally termini yielded mono-, di-, and tri-epoxy and -hydroxy-telechelic polyisobutylenes. It is further disclosed that strong rubbery networks were made by curing the epoxy-telechelic polyisobutylenes with triethylene tetramine and by reacting the hydroxy-telechelic polyisobutylenes with MDI.\nJ. P. Kennedy, Sandor Nemes and Tibor Pernecker in Polymer Bulletin, V. 25, p.633 (1991) disclose vinyl ether headed polyisobutylenes macromonomers. However no mention was made regarding radiation curable compositions based on these macromonomers. It is known that radiation cured networks from non-telechelic macromonomers will posses poor physical properties.\nN. A. Merrill, I. J. Gardner, and V. L. Hughes in RadTech North America Proceedings, V. 1, pp. 77-85 (1992) disclose conjugated diene functional polyisobutylenes which have a high reactivity to both ultraviolet and electron beam radiation. These conjugated diene functional polyisobutylenes, alone or in a formulation, are disclosed as being useful in preparing pressure sensitive adhesives.\nIn PCT Patent Publication No. WO 9104992 is disclosed a functionalized copolymer of isobutylene and a para-methylstyrene, wherein at least one type of functional group is attached to the para-methyl group of the para-methylstyrene, the copolymer having a substantially homogenous compositional distribution. The functionalized groups are exemplified by alkoxides, phenoxides, carboxylates, thiolates, thiopenolates, thioethers, thiocarboxylates, dithiocarboxylates, thioureas, dithiocarbamates, xanthanates, thiocyanates, silanes, halosilanes, malonates, cyanides, amides, amines, carbazoles, phthalimides, pyridine, maleimide, cyanates, and phosphines.\nIn PCT Patent Publication No. WO 9211295 is disclosed a radiation reactive functionalized polymer comprising an isoolefin having about 4 to about 7 carbon atoms and a para-alkylstyrene, wherein a radiation reactive functional group is attached to the para-alkyl group of the para-alkylstyrene, and discloses radiation curable pressure sensitive adhesives comprising the functionalized polymer and a tackifier. In WO\"\"295, the radiation curable groups are disclosed as being groups such as thioxanthones, acrylates, aldehydes, ketones, and esters.\nSaxena et al. in U.S. Pat. No. 5,665,823 disclose a method for preparing an acrylic functional polyisobutylene polymer or copolymer, the method comprising reacting a polyisobutylene polymer or copolymer which contains at least one carbon-bonded silanol group in it molecule with a silane having both an acrylic-containing group and a silicon-bonded hydrolyzable group in its molecule.\nFurthermore, radiation curable compositions which contain vinyl ether functional organosilicon compounds have also been described in the art. For example, Crivello in U.S. Pat. No. 4,617,238 discloses a photopolymerizable composition comprising (a) an organopolysiloxane having at least one Si-bonded vinyloxy functional group of the formula H2Cxe2x95x90CHxe2x80x94Oxe2x80x94Gxe2x80x94, where G is alkylene (such as propylene) or alkylene interrupted by at least one divalent heteroradical selected from xe2x80x94Oxe2x80x94, divalent phenylene, or substituted divalent phenylene, or combination of such heteroradicals, and (b) an onium salt catalyst. The \"\"238 patent also describes a method wherein the vinyl ether group is introduced into the organopolysiloxane by addition (hydrosilylation) of compounds with an allyl and a vinyl ether group to an SiH group of the organopolysiloxane in the presence of a platinum catalyst. In the method of the \"\"238 patent, only the allyl group is added to the SiH group while the vinyl ether group is preserved and thus only one vinyl ether group for each SiH group can be incorporated into the siloxane molecule at any given time.\nEuropean Patent Publication No. 0462389 teaches thermosetting organopolysiloxanes with oxyalkylene vinyl ether groups bonded by SiOC groups and the vinyl groups may be substituted by alkyl groups. EPO\"\"389 also teaches a method for the preparation of these compounds and their application as photochemically thermosetting polysiloxanes in encapsulating compounds, as non-stick coating compounds for flat carriers or as modified additives in compounds which can be thermoset radically, cationically or by UV or electron radiation.\nBrown et al., in U.S. Pat. No. 5,270,423 disclose organosilicon compounds with a siloxane portion of the general formula xe2x80x94ORxe2x80x2OCHxe2x95x90CHRxe2x80x3 linked via an SiOC bond wherein Rxe2x80x2 is a divalent hydrocarbon group and Rxe2x80x3 is hydrogen or an alkyl group which are useful in radiation curable compositions, in which they are mixed with an initiator. The compositions are particularly useful in UV radiation curable coatings.\nGlover et al. in U.S. Pat. No. 5,594,042 discloses radiation curable compositions comprising vinyl ether functional siloxanes and aromatic iodonium salt or aromatic sulfonium salt photoinitiators which cure upon exposure to ultraviolet or electron beam radiation. The vinyl ether groups are linked to the silicon atom on the siloxane through an SiOC bond and the photoinitiators are disclosed as being preferably either diaryliodonium salts of sulfonic acids or triarylsulfonium salts of sulfonic acids.\nBujanowski et al. in U.S. Pat. No. 5,629,095 disclose vinyl ether functional siloxane resins, radiation curable coating compositions comprising a vinyl ether functional siloxane resin and a photocleavable acid, and a coated article obtained by applying the radiation curable coating composition to a substrate and then exposing the coating to radiation in an amount sufficient to cure the coating. In the \"\"095 patent, the vinyl ether group in the siloxane resin is attached to the silicone atom through an SiOC bond.\nThe present invention relates to radiation curable compositions comprising an alkenyl ether functional polyisobutylene, a cationic photoinitiator and a miscible reactive diluent. The compositions of this invention can further comprise a free radical photoinitiator, photosensitizers, non-reactive diluents for improving the cure speed of the radiation curable compositions, thermo-oxidative stabilizers to improve thermal and oxidative stability of the compositions, and stabilizers for improving the shelf life of the compositions.\nIt is an object of this invention to produce radiation curable compositions which exhibit a low cure energy.\nIt is an object of this invention to produce radiation curable compositions which have a high moisture vapor barrier, high damping characteristics, and a high refractive index. It is an object of this invention to produce radiation curable compositions which provide a barrier to corrosive vapors and have maintained or enhanced modulus, tensile strength, and toughness.\nThis invention relates to a radiation curable composition comprising: (A) an alkenyl ether-functional polyisobutylene polymer in which at least 50 mole percent of the non-terminal repeating units of the polymer are isobutylene units and containing at least one group having the formula \nwherein R is independently selected from the group consisting of monevalent hydrocarbon groups and alkoxy groups, R1 is a divalent hydrocarbon group having from 2 to 20 carbon atoms, R2 and R3 are independently selected from the group consisting of a hydrogen atom and a monovalent hydrocarbon group, and a has a value of 0 to 2, (B) a cationic photoinitiator, and (C) a miscible reactive diluent selected from\n(i) a difunctional vinyl ether reactive diluent selected from the group consisting of butanediol divinyl ether, pentanediol divinyl ether, hexanediol divinyl ether, heptanediol divinyl ether, cyclohexane dimethanol divinyl ether, ethylene glycol divinyl ether, diethylene glycol divinyl ether, triethylene glycol divinyl ether, neopentyl glycol divinyl ether, ethoxylated2-20 bisphenol A divinyl ether, poly-THF divinyl ether, bis-(4-ethenyloxy butyl)-succinate, bis-(4-ethenyloxy butyl)-adipate, bis-(4-ethenyloxy butyl)glutarate, bis-((4-((ethenyloxy)methyl)cyclohexyl)methyl)succinate, bis((4-((ethenyloxy)methyl)cyclohexyl)methyl)adipate, and bis-((4-((ethenyloxy)methyl)cyclohexyl)methyl)glutarate; or\n(ii) an acrylate reactive diluent with the formula\n(R8)Xb\nxe2x80x83wherein R8 is a non-silicon containing organic group, X is an organic group containing at-least one acrylate functional group, and b has a value of 2 to 4.\nIn Component (A), the alkenyl ether-functional polyisobutylene polymer, the monovalent hydrocarbon groups of R are exemplified by alkyl groups such as methyl, ethyl, propyl, butyl, hexyl, octyl, and decyl, aryl groups such as phenyl, tolyl, and xylyl, and can also be any monovalent hydrocarbon group which has at least one of its hydrogen atoms replaced with a halogen, such as fluorine, chlorine, or bromine, and these monovalent hydrocarbon groups are exemplified by CF3CH2CH2xe2x80x94 and C4F9CH2CH2xe2x80x94. The alkoxy groups are exemplified by methoxy, ethoxy, propoxy, and butoxy. It is highly preferred that R is independently selected from the group consisting of methyl and methoxy. Each R group can be the same or different, as desired.\nDivalent hydrocarbon groups suitable as R1 are exemplified by alkylene groups such as ethylene, propylene, butylene, pentylene, trimethylene, 2-methyltrimethylene, pentamethylene, hexamethylene, 3-ethyl-hexamethylene, octamethylene, decamethylene, xe2x80x94(CH2)18xe2x80x94, and cycloalkylene groups such as cyclohexylene, arylene groups such as phenylene. Examples of suitable divalent halohydrocarbon groups also include any divalent hydrocarbon group wherein one or more hydrogen atoms have been replaced by halogen, such as fluorine, chlorine or bromine exemplified by xe2x80x94CH2CH2CF2CF2CH2CH2xe2x80x94. Each R1 can be the same or different as desired. Preferably R1 is butylene.\nThe groups R2 and R3 are independently selected from a group consisting of a hydrogen atom and a monovalent hydrocarbon group exemplified by alkyl groups such as methyl, ethyl, propyl, butyl. The groups R2 and R3 may be the same or different. Preferably R2 and R3 are hydrogen atoms. In the formula above, it is preferred that a have a value of zero.\nFor the purposes of this invention, the backbone of the alkenyl ether functional polyisobutylene polymer may be any linear or branched polymer or copolymer wherein at least about 50 mole percent, preferably at least 80 mole percent, of the repeat units are isobutylene units having the following structure: \nIn the polymer or copolymer of the invention, the above described alkenyl ether group can be disposed either along the chain or at the terminals thereof, or any combination of the above. As used herein, the term xe2x80x9cpolymerxe2x80x9d is generic to polymers, oligomers, and copolymers, all of which are within the scope of this invention.\nIn a preferred embodiment of this invention, the alkenyl-ether functional polyisobutylene polymer (A) is a polymer containing at least one group having the formula \nwherein at least 50 mole percent of the non-terminal repeating units of the polymer are isobutylene units, R is independently selected from the group consisting of monovalent hydrocarbon groups and alkoxy groups, R1 is a divalent hydrocarbon group having from 2 to 20 carbon atoms, R2 and R3 are independently selected from group consisting of a hydrogen atom and a monovalent hydrocarbon group, n has a value from 5 to 10,000, a has a value of 0 to 2, and Y is selected from the group consisting of (i) an alkylene group having from 2 to 10 carbon atoms and (ii) a group having the formula \nwherein R4 is a monovalent hydrocarbon group, R5 and R6 are independently alkylene groups having from 2to 10 carbon atoms, and m is an integer having a value from 1 to 5.\nThe groups R and R1 are as defined hereinabove, including preferred embodiments thereof. Preferably, R is independently selected from the group consisting of methyl and methoxy, and R1 is butylene. Preferably, a has a value of 0 or 1.\nThe alkylene groups of Y(i) are exemplified by ethylene, propylene, butylene, pentylene, trimethylene, 2-methyltrimethylene, pentamethylene, hexamethylene, 3-ethyl-hexamethylene, octamethylene, and decamethylene.\nIn the formula for Y(ii) above, the monovalent hydrocarbon groups of R4 are as described above for R, and preferably R4 is methyl. The alkylene groups for R5 and R6 are as defined above for Y(i). Preferably, R5 and R6 are independently selected from the group consisting of ethylene and propylene. It is highly preferred that R5 is propylene, and R6 is ethylene. It is also preferred that m has a value of 1.\nIt is preferred for purposes of this invention that from 10 to 100 weight percent of the alkenyl ether functional polyisobutylene polymer described above be used, and it is highly preferred that from 50 to 100 weight percent of this compound be employed, said weight percent being based on the total weight of the radiation curable composition.\nComponent (B) in the compositions of this invention is cationic photoinitiator. Suitable cationic photoinitiators are selected from the group consisting of onium salts, diaryliodonium salts of sulfonic acids, triarylsulfonium salts of sulfonic acids, diaryliodonium salts of boronic acids, and triarylsulfonium salts of boronic acids.\nThe onium salts are preferably selected from the group consisting of R72I+MXZxe2x88x92, R73S+MXZxe2x88x92, R73Se+MXZxe2x88x92, R74P+MXZxe2x88x92, and R74N+MXZxe2x88x92, wherein each R7 is an organic group having from 1 to 30 carbon atoms exemplified by aromatic carbocyclic groups having from 6 to 20 carbon atoms. Each R7 can be substituted with from 1 to 4 monovalent hydrocarbon groups exemplified by alkoxy groups having from 1 to 8 carbon atoms, alkyl groups having from 1 to 16 carbon atoms, nitro, chloro, bromo, cyano, carboxyl, mercapto, and aromatic heterocyclic groups exemplified by pyridyl, thiophenyl, and pyranyl. The symbol M in the formulae hereinabove are metals or metalloids which include transition metals exemplified by Sb, Fe, Sn, Bi, Al, Ga, In, Ti, Zr, Sc, V, Cr, Mn, Cs, rare earth metals exemplified by lanthanides, for example, Cd, Pr, and Nd, and metalloids exemplified by B, P, and As. MXZxe2x88x92is a non-basic, non-nucleophilic anion exemplified by BF4xe2x88x92, PF6xe2x88x92, AsF6xe2x88x92, SbF6xe2x88x92, SbCl6xe2x88x92, HSO4xe2x88x92, ClO4xe2x88x92, FeCl4xe2x95x90, SnCl6xe2x88x92, and BiCl5xe2x95x90.\nPreferred onium salts are exemplified by bis-diaryl iodonium salts, for example, bis(dodecyl phenyl)iodonium hexafluoroarsenate, bis(dodecylphenyl)iodonium hexafluoroantimonate, and dialkylphenyl iodonium hexafluoroantimonate.\nDiaryliodonium salts of sulfonic acids, triarylsulfonium salts of sulfonic acids, diaryliodonium salts of boronic acids, and triarylsulfonium salts of boronic acids are also suitable as the cationic photoinitiator (B). Preferred diaryliodonium salts of sulfonic acid are diaryliodonium salts of perfluoroalkylsulfonic acids and diaryliodonium salts of aryl sulfonic acids. Preferred diaryliodonium salts of perfluoroalkylsulfonic acids are exemplified by diaryliodonium salts of perfluorobutanesulfonic acid, diaryliodonium salts of perfluoroethanesulfonic acid, diaryliodonium salts of perfluoro-octanesulfonic acid, and diaryliodonium salts of trifluoromethane sulfonic acid. Preferred diaryliodonium salts of aryl sulfonic acids are exemplified by diaryliodonium salts of para-toluene sulfonic acid, diaryliodonium salts of dodecylbenzene sulfonic acid, diaryliodonium salts of benzene sulfonic acid, and diaryliodonium salts of 3-nitrobenzene sulfonic acid.\nPreferred triarylsulfonium salts of sulfonic acid are triarylsulfonium salts of perfluoroalkylsulfonic acids and triarylsulfonium salts of aryl sulfonic acids. Preferred triarylsulfonium salts of perfluoroalkylsulfonic acids are exemplified by triarylsulfonium salts of perfluorobutanesulfonic acid, triarylsulfonium salts of perfluoroethanesulfonic acid, triarylsulfonium salts of perfluoro-octanesulfonic acid, and triarylsulfonium salts of trifluoromethane sulfonic acid. Preferred triarylsulfonium salts of aryl sulfonic acids are exemplified by triarylsulfonium salts of para-toluene sulfonic acid, triarylsulfonium salts of dodecylbenzene sulfonic acid, triarylsulfonium salts of benzene sulfonic acid, and triarylsulfonium salts of 3-nitrobenzene sulfonic acid.\nPreferred diaryliodonium salts of boronic acids, and triarylsulfonium salts of boronic acids are compounds such as those disclosed in European Patent Application No. 0562922. Preferred diaryliodonium salts of boronic acids include diaryliodonium salts of perhaloarylboronic acids and preferred triarylsulfonium salts of boronic acids are the triarylsulfonium salts of perhaloarylboronic acid.\nPreferably the amount of cationic photoinitiator (B) is from 0.01 to 5.0 weight percent based on the total weight of the composition, and it is highly preferred to use from 0.1 to 2.0 weight percent based on the total weight of the radiation curable composition.\nPreferably, Component (C)(i) is selected from the group consisting of butanediol divinyl ether, hexanediol divinyl ether, and cyclohexane dimethanol divinyl ether.\nComponent (C)(ii) in the compositions of this invention is a miscible reactive diluent having the formula\n(R8)Xb\nwherein R8=a non-silicon containing organic group; X=organic group containing at-least one acrylate functional group; and b=2-4. The R8 group is selected to be the one that is compatible with component (A). The R8 group has between 2 to 30 carbon atoms, and preferably between 4 to 18 carbon atoms. The acrylate functional reactive diluent can be monofunctional, difunctional, or trifunctional.\nComponent (C)(ii) is exemplified by ethyl acrylate, propyl acrylate, butyl acrylate, tert.-butyl acrylate, tert.-amyl acrylate, hexyl acrylate, cyclohexyl acrylate, 2-ethylhexyl acrylate, isooctyl acrylate, isodecyl acrylate, dodecyl acrylate, lauryl acrylate, stearyl acrylate, ethyleneglycol butyl acrylate, tetrahydrofurfuryl acrylate, isobornyl acrylate, tridecyl acrylate, caprolactone acrylate, 2-phenoxyethyl acrylate, ethoxylated nonyl phenol acrylate, butanediol diacrylate, hexanediol diacrylate, cyclohexane dimethanol diacrylate, ethylene glycol diacrylate, diethylene glycol diacrylate, triethylene glycol diacrylate, neopentyl glycol diacrylate, ethoxylated2-20 bisphenol A diacrylate, poly-THF diacrylate, and trimethylolpropane tri acrylate.\nPreferably the amount of miscible reactive diluent (C) is from 1 to 60 weight percent based on the total weight of the composition, and it is highly preferred to use from 5 to 40 weight percent based on the total weight of the radiation curable composition.\nVarious optional components may be used in the composition beyond Components (A) through (C). The photocuring of alkenyl-ethers is a cationic process. One way known in the art to increase cure speeds of cationic systems is to use one or more free radical photoinitiators and photosensitizers. The use of non-reactive diluents may be required to efficiently disperse the photocatalyst, photoinitiators and photosensitizers into the polyisobutylene matrix.\nThe free radical photoinitiators for use in the compositions of this invention can be any benzoins exemplified by benzoin alkyl ethers, benzophenone and its derivatives such as 4,4xe2x80x2-dimethyl-amino-benzophenone (Michler\"\"s Ketone), acetophenones exemplified by dialkoxyacetophenones, dichloroacetophenones, and trichloroacetophenones, benzils exemplfied by benzil ketals, quinones, and O-acylated -xcex1-oximinoketones. Preferably the free radical photoinitiator is a compound having the formula \nwherein Rxe2x80x2 is selected from the group consisting of xe2x80x94H, an alkoxy group, and a halogen atom, Rxe2x80x3 is selected from the group consisting of xe2x80x94OH, an alkoxy group, and a halogen atom, and Rxe2x80x2xe2x80x3 is selected from the group consisting of xe2x80x94H, an alkyl group, and a halogen atom. Preferred embodiments of this compound are (i) where Rxe2x80x2 is xe2x80x94H, Rxe2x80x3 is xe2x80x94OH and Rxe2x80x2xe2x80x3 is methyl or phenyl, (ii) where Rxe2x80x2 is xe2x80x94H, Rxe2x80x3 is an alkoxy group and Rxe2x80x2xe2x80x3 is phenyl (for benzoin alkyl ethers), (iii) where both Rxe2x80x2 and Rxe2x80x3 are alkoxy groups and Rxe2x80x2xe2x80x3 is phenyl (for benzil ketals), (iv) where both Rxe2x80x2 and Rxe2x80x3 are alkoxy groups and Rxe2x80x2xe2x80x3 is xe2x80x94H (for dialkoxyacetophenones), and (v) where both Rxe2x80x2 and Rxe2x80x3 are xe2x80x94Cl and Rxe2x80x2xe2x80x3 is xe2x80x94Cl or xe2x80x94H (for di- and tri-chloroacetophenones). A preferred free radical photoinitiator is Darocure(copyright) 1173 (2-hydroxy-2-methyl-1-phenyl-propan-1-one).\nPreferably the amount of free radical photoinitiator is from 0.01 to 5.0 weight percent based on the total weight of the composition, and it is highly preferred to use from 0.1 to 2.0 weight percent based on the total weight of the radiation curable composition.\nPreferred photosensitizers are 2-isopropylthioxanthone or benzophenone.\nPreferably the amount of photosensitizer is from 0.01 to 2.0 weight percent based on the total weight of the composition, and it is highly preferred to use from 0.05 to 0.5 weight percent based on the total weight of the radiation curable composition.\nPreferred non-reactive diluent are hydrocarbon diluents with C4-20 hydrocarbon chain length; long chain hydrocarbon diluents with epoxy, ester, ether, glycidyl ether, anhydride and carbonyl functional groups. Suitable examples include methyl laurate, methyl nonate, ethyl laurate, dioctyl adipate, di-(2-ethylhexyl)phthalate. di-2-ethylhexyl ether, dioctadecyl ether, dodecylene epoxide, hexyl glycidyl ether, and succinic anhydride.\nPreferably the amount of non-reactive diluent is from 0.01 to 2.0 weight percent based on the total weight of the composition, and it is highly preferred to use from 0.05 to 0.5 weight percent based on the total weight of the radiation curable composition.\nTo improve thermal and oxidative stability of the cured coating one or more stabilizers may be included in the composition. Examples of suitable stabilizers are hindered amines, organic phosphites, hindered phenols and mixtures thereof. Some of the preferred stabilizers are: Irganox(trademark) 1520D; Irganox(trademark) 1010, Tinuvin(trademark) 123 and Tinuvin(trademark) 292.\nPreferably the amount of stabilizer is from 0.01 to 2.0 weight percent based on the total weight of the composition, and it is highly preferred to use from 0.05 to 0.5 weight percent based on the total weight of the radiation curable composition.\nTo improve shelf stability of the un-cured coating one or more stabilizers may be included in the composition. This requires the use of stable photoinitiator compositions as described in the U.S. Pat. No. 5,973,020. These stabilizing agents are usually hindered amines with boiling points of greater than 150xc2x0 C. and preferably greater than 200xc2x0 C. The amine can be a primary, secondary or tertiary amine and preferably a secondary and tertiary amine. Examples are Tinuvin products sold by Ciba-Geigy company, the CYAGARD products sold commercially by Cytec Corporation (Stamford, Conn.) and SANDUVAR products as referenced in the U.S. Pat. No. 5,973,020.\nPreferably the amount of stabilizer generally used, by weight in relation to the total weight of the cationic photoinitiator (B) is from 0.01 to 2.0 percent by weight of the cationic photoinitiator solids and preferably between 0.2 to 1.0 weight percent.\nThe radiation curable compositions of this invention can also contain ingredients exemplified by reinforcing and extending fillers such as treated silicas, hydrocarbon diluents such as linear alkyl dodecylbenzene and functional hydrocarbons such as C8-16 aliphatic glycidyl ethers, colorants, dyes, preservatives, fragrances, and adhesion modifiers.\nThe radiation curable compositions of this invention can be prepared by mixing the materials described hereinabove and any optional components in any order, using any suitable mixing means, such as a spatula, a drum roller, a mechanical stirrer, a three-roll mill, a sigma blade mixer, a bread dough mixer, or a two-roll mill.\nThis invention further relates to a method of making a radiation curable composition comprising (I) mixing components (A)-(C) and the optional ingredients described hereinabove. Components (A)-(C) and the optional ingredients are as described above, including preferred embodiments and amounts thereof.\nThe present invention further relates to a method of making an article of manufacture comprising (I) applying a radiation curable composition comprising components (A)-(C) described hereinabove, to a solid substrate to form a coating, and (II) exposing the coating to an energy source selected from the group consisting of (i) ultraviolet light and (ii) visible light in an amount sufficient to cure the coating.\nThe composition of (I) can further comprise any of the optional ingredients recited above. Components (A)-(C) and the optional ingredients are as described above, including preferred embodiments and amounts thereof.\nThe coating may be applied by any suitable manner known in the art, such as by spreading, brushing, extruding, spraying, gravure, kiss-roll and air-knife.\nThe solid substrate can be a flexible sheet material such as paper, polyolefin film, polyolefin-coated paper, foil, wood, cardboard and cotton, metallic materials such as aluminum, copper, steel and silver, siliceous materials such as glass and stone, and synthetic polymer materials such as polyolefins, polyamides, polyesters and polyacrylates. As to form, the solid substrate can be substantially sheet-like, such as a peelable release liner for pressure sensitive adhesive, a fabric or a foil, or a fiber, or a substantially three-dimensional in form.\nCuring itself may be achieved in any of the known ways, including passing a coated substrate under the desired source of radiation, for example a UV lamp, at a predetermined rate and exposing a completely coated substrate to radiation by switching on the required energy source for a predetermined time.\nThe radiation curable compositions are preferably cured in the form of films. The cured films are expected to have high refractive index, good barrier properties, good adhesion and good damping properties. It is preferable to apply these coatings to surfaces that are adversely affected by exposure to oxygen, moisture vapor and other environmental factors. The radiation curable coatings are particularly useful as high refractive index coatings for optical fibers. The application of the radiation curable compositions to optical fibers and curing of the compositions can be achieved by conventional equipment (see Blyler and Aloisio Polymers for Coating Optical Fibers, Chemtech, November 1987, pages 680-684). The curable compositions can also be used as an additive to compositions whose barrier properties needs to be tailored to higher values. The radiation curable compositions can be used to increase the barrier properties of sealants and pottants used for encapsulating electronic devices that are adversely affected by moisture."}
{"text": "Conventional golf balls can be divided into two general classes: solid and wound. Solid golf balls include one-piece, two-piece (i.e., single layer core and single layer cover), and multi-layer (i.e., solid core of one or more layers and/or a cover of one or more layers) golf balls. Wound golf balls typically include a solid, hollow, or fluid-filled center, surrounded by a tensioned elastomeric material, and a cover.\nGolf ball core and cover layers are typically constructed with polymer compositions including, for example, polybutadiene rubber, polyurethanes, polyamides, ionomers, and blends thereof. Ionomers, particularly ethylene-based ionomers, are a preferred group of polymers for golf ball layers because of their toughness, durability, and wide range of hardness values.\nGolf ball compositions comprising highly neutralized acid polymers are known. For example, U.S. Patent Application Publication No. 2003/0130434, the entire disclosure of which is hereby incorporated herein by reference, discloses melt-processible, highly-neutralized ethylene acid copolymers and process for making them by incorporating an aliphatic, mono-functional organic acid in the acid copolymer and then neutralizing greater than 90% of all the acid groups present. The use of such compositions in various golf ball layers is disclosed. Also, U.S. Patent Application Publication No. 2005/0148725, the entire disclosure of which is hereby incorporated herein by reference, discloses a highly-resilient thermoplastic composition comprising (a) an acid copolymer, (b) a salt of a high molecular weight, monomeric organic acid; (c) a thermoplastic resin; (d) a cation source; and (e) optionally, a filler. The reference also discloses one-piece, two-piece, three-piece, and multi-layered golf balls comprising the highly-resilient thermoplastic composition.\nWhile various uses for highly neutralized acid polymers in golf balls have been discovered, there is a need in the industry to broaden the applicability of highly neutralized acid polymers to particular golf ball constructions having desirable spin, feel, and COR properties. The present invention provides such golf ball constructions through the use of a layer formed from a low modulus HNP composition and a layer formed from a high modulus HNP composition."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory and, more specifically, to techniques particularly effectively applied to pseudostatic RAMs (random access memories).\n2. Description of the Related Art\nA known pseudostatic RAM is designed on the basis of a dynamic RAM which can be constructed in a large scale of integration, and has an interface compatibility with an ordinary static RAM. The pseudostatic RAM has, in addition to an ordinary read/write mode, an address refresh mode and an autorefresh mode in which a refreshing action is executed once in response to an external control action, and a self-refresh mode in which a self-controlled refreshing action is executed periodically, for example, while the pseudostatic RAM is backed up by a battery. The pseudostatic RAM is provided internally with a refresh counter for sequentially specifying word lines for refreshing action in the autorefresh mode and the self-refresh mode, and a refresh timer circuit for periodically starting the refreshing action in the self-refresh mode.\nThe pseudostatic RAM having the autorefresh and self-refresh modes is described, for example, in \"Hitachi IC Memory Data Book\", Hitachi, Ltd., pp. 229-234, March, 1987.\nThe mean current consumption of the memory arrays of the pseudostatic RAM or the like in the self-refresh mode increases approximately in proportion to the reciprocal of the refresh period, i.e., the number of refresh cycles in a unit time. The refresh period in the self-refresh mode is dependent on the intrinsic information holding ability of the memory cells and the stability of the refresh timer circuit for setting the refresh period. Such conditions are restrictions on the reduction of the power consumption of the pseudostatic RAM or the like in operation backed up by batteries.\nPrior to the present invention, the inventors of the present invention intended to stabilize the refresh period of the pseudostatic RAM and to increase the refresh period as near to the intrinsic information holding ability of the memory cells as possible by constructing the refresh timer circuit of an oscillation circuit, which is less dependent on the supply voltage, and a refresh timer counter circuit which counts the pulses of a pulse signal provided by the oscillation circuit and generates a predetermined refresh start signal, and by selectively opening corresponding fuse means to set an optional initial count for the refresh timer counter circuit.\nHowever, the operating current of the oscillation circuit is limited, and the oscillation circuit needs a capacitor having a comparatively long charging or discharging period, and a comparatively large resistance, hence a comparatively long polycrystalline silicon resistor must necessarily be formed on a semiconductor substrate. Accordingly, the discharge current varies or the bumping (varying) of the supply voltage cannot be absorbed rapidly due to a large substrate capacitance, which is parasitic to the semiconductor substrate and the polycrystalline silicon resistor causing the variation of the frequency of oscillation of the oscillation circuit, for example, when the supply voltage applied to the oscillation circuit bumps (varies) during the discharge period of the capacitor. Therefore, the frequency of oscillation of the oscillation circuit must have a margin in relation with the information holding ability of the memory cells.\nAlthough the oscillation circuit and refresh timer counter circuit of the pseudostatic RAM is capable of selectively changing the refresh period of the pseudostatic RAM, there is no means for testing and confirming the oscillation characteristics and variation characteristics. Accordingly, the refresh period is set inevitably by a trial-and-error method unless data of those characteristics is available. This fact, similarly, requires the frequency to be set with a margin in relation with the information holding ability of the memory cells, restricts the reduction of the power consumption of the pseudostatic RAM and increases the amount of work required for testing the pseudostatic RAM."}
{"text": "The ability of a polished surface to return to its original appearance after being smeared and the amount of time required to do so are two factors which tend to be rather important in the evaluation of the performance of a furniture polish composition.\nUpon application onto a particular substrate surface, the polish compositions of the present invention have been observed to form a rather thin (approximately 1000 to 4000 Angstroms) generally uniform film on the substrate surface. Due to the liquid nature of the resultant films of the surface-applied emulsion polish compositions of the present invention, such films typically exhibit, to a degree, visible physical evidence of surface disturbances, otherwise referred to as surface \"smears\", when the surface-applied film is subjected to a rubbing action.\nThe term \"smear recovery\", as used herein, denotes not only the rate-of-recovery but also the ability of a thus \"smeared\" surface-applied polish film to return to its original surface appearance.\nPolish compositions possessing superior smear recovery properties have long been desired. Conventional polishes, however and in particular the surface-applied films they produce lack this quality. For example, commercial polishes that use conventional silicone-containing or conventional silicone-based fluids, as well as certain conventional waxes, and certain low molecular weight organic oils as glossing agents seemingly initially tend to provide certain desirable surface-appearance values but, in fact, ultimately suffer from certain inherent disadvantages. Such conventional polishes, in particular, readily tend to show certain surface disturbances such as fingerprints and otherwise readily tend to smear.\nCertain investigations discussed in the prior-art references reviewed by us have attacked these and other surface \"smear\" problems by attempting to produce specific smear-resistant polishes which are said to provide treated surfaces that have relatively permanent finishes. Such so-called \"permanent\" finishes are said to require the application of significant force before their surface properties are disturbed. While such finishes tend to provide both high gloss and smear resistant properties, it is generally more difficult to remove these so-called \"permanent\" finishes; and such finishes thus have the tendency to build up with multiple applications.\nSmear resistant films of this type are described e.g. in U.S Pat. No. 3,847,622 to Brandl et al. The disclosed compositions of Brandl et al. contain water, wax, a mixture of organopolysiloxanes a solvent, and a water-in-oil emulsifier. Such compositions furthermore, are said to be substantially devoid of any tendency to smear or streak and are further said to be easily applied with a minimum of physical effort. Brandl et al. teach however that the use of silicones of diverse types, i.e. organopolysiloxanes, while beneficial from the standpoint of luster, gloss and durability nevertheless fail to provide appreciable mitigation of the smear problem as they view it. The emulsion polish compositions disclosed by Brandl et al. are thus said to possess a \"particularly high resistance to smearing.\"\nIn U.S. Pat. No. 3,856,533, Schnurrbusch et al. disclose certain wax-containing surface-polish compositions which are said to be resistant to \"finger-marking\". Such compositions are disclosed as containing an organopolysiloxane (such as dimethylsiloxane), a wax, and a specified so-called \"organosiloxane\" containing an alkyl, an alkenyl, or a monocyclic aryl radical.\nThe present invention, in contradistinction to the prior art briefly discussed hereinabove, approaches the so-called \"smear problem\" in a manner which is markedly different from the approach set forth either in U.S. Pat. No. 3,847,622 to Brandl et al. or in U.S. Pat. No. 3,856,533 to Schnurrbusch et al. The present invention, in particular, enhances the smear-recovery properties of the emulsion polish compositions disclosed herein. Indeed, the unexpected and significant increase in smear recovery which is attributable to the polysiloxane copolymers of our present invention, is shown in the comparative-example data presented hereinbelow. In particular, the smear-recovery properties of the novel emulsion polish compositions disclosed herein were not recognized by any of the prior-art polish formulators whose work we reviewed\nThe prior-art references which we reviewed while disclosing numerous furniture polishes using polysiloxanes, thus do not teach, disclose or even suggest the novel polysiloxane copolymers of our present invention. In particular, and as mentioned above, the novel polysiloxane copolymers of our invention have been demonstrated, by comparative example, as being unexpectedly superior over prior-art polysiloxanes such as polydimethylsiloxane. Indeed, it was only through our present discovery that the surprising smear-recovery properties of the emulsion polish compositions of the present invention came to light.\nIn U.S. Pat. No. 2,698,805 to Currie et al., for example, there are disclosed certain conventional polishes which are said to be easy to apply. Such conventional emulsion polish compositions contain a hydrocarbon-soluble organopolysiloxane having from 1 to 3 aryl, alkyl, or arylkyl radicals per silicon atom. Such emulsion polish compositions further include a hydrocarbon solvent a specified aluminum stearate compound, and water. The alkyl radicals which Currie et al. disclose are methyl, ethyl, propyl and octadecyl. The polish compositions disclosed by Currie et al., furthermore, are said to be smear-free. In contradistinction, the polish compositions of our present invention do indeed smear; but, upon being smeared, recover more rapidly than conventional compositions containing conventional polydimethylsiloxane polymers of comparable viscosity and/or weight-average molecular weight relative to the above-mentioned polymeric siloxanes disclosed herein.\nIndeed the concept of rapid smear recovery is not disclosed or even suggested by Currie et al. Moreover, the preferred organopolysiloxane disclosed in the Currie '805 patent, namely dimethylpolysiloxane fluid, is shown hereinbelow as being inferior to the siloxane-containing compositions of the present invention with regard to smear-recovery properties.\nFurther, U.S. Pat. No. 2,523,281 to Currie, is directed to automobile polishes containing an organopolysiloxane having from 1 to 3 aryl, alkyl or alkaryl radicals per silicon atom. Such automobile polishes also include a hydrocarbon solvent, finely-divided silica of a specified particle size, an emulsifying agent, and water. The disclosed alkyl radicals include methyl, ethyl, propyl, and octadecyl. Disclosed in this patent, furthermore, as the preferred organo polysiloxanes, are certain dimethylpolysiloxanes which are shown hereinbelow as possessing inferior smear-recovery properties as compared to those siloxanes utilized in accordance with the principles of the present invention.\nStill other patents, disclosing the use of polysiloxanes in furniture-care compositions, do not disclose the novel emulsion polish compositions of the present invention. In particular, U.S. Pat. Nos. 4.163.673 to Dechert; 4,354,871 to Sutton; and U.S. Pat. No. 4,613,646 to Sandvick all disclose that dimethyl polysiloxane or polydimethyl siloxane is an especially-preferred polysiloxane for use in furniture polish compositions. Such siloxanes however form no part of the novel emulsion polish compositions of the present invention.\nIndeed, the prior art does not disclose or even suggest the specific poly(dimethyl)-co-poly(methylalkyl) siloxane and poly(dimethyl)-co-poly(methyl, oxygen-containing) siloxane copolymers of the present invention; nor does the prior art disclose or even suggest the unexpected smear-recovery properties of compositions containing the siloxane copolymers of the present invention.\nFor example, U.S. Pat. No. 4,404,035 to Ona et al. discloses wax compositions comprising certain hydrocarbon ester-containing organopolysiloxanes, i.e. specified organopolysiloxanes containing at least one silicon-bonded hydrocarbon-ester group. The Ona et al. patent does not, however, teach or even suggest the usefulness of alkylene oxide polysiloxanes. Ona et al., in particular, teach away from using alkylene oxide polysiloxanes, as is shown by way of their so-called \"solution D\", as presented in their Example 2. That is, Ona et al., in their Example 2, specifically report that a polysiloxane having a side chain containing an alkylene oxide having the following formula EQU (CH.sub.2).sub.3 O(C.sub.2 H.sub.4 O).sub.20 (C.sub.3 H.sub.6 O).sub.20 H\nproduced a coating on a metal panel which, after being subjected to a 20-hour water shower, was deemed to be unacceptable in luster, water-repellency and durability.\nIn U.S. Pat. No. 4,218,250 to Kasprzak, there is disclosed a wax-containing and/or abrasive-containing formulation which includes a cyclodimethyl siloxane fluid and a polydiorgano siloxane polyoxyalkylene copolymer. The disclosed polish formulation may also include a specific silicon-glycol copolymer having the following formula EQU (CH.sub.3).sub.3 SiO[(CH.sub.3).sub.2 SiO].sub.x [G(CH.sub.3)SiO].sub.y Si(CH.sub.3).sub.3\nwherein x has an average value of 0 to 10, y has an average value of 1 to 10, G has the formula, --D--(OC.sub.2 H.sub.4).sub.z OH, with D being an alkylene radical containing from 2 to 10 carbon atoms and z having an average value of from 6 to 20. This patent, however, does not mention smear resistance or smear recovery; nor does this patent otherwise disclose either the specific poly(dimethyl)-co-poly(methylalkyl) siloxane or poly(dimethyl)-co-poly(methyl, oxygen-containing) siloxane copolymers of the present invention.\nAdditionally U.S. Pat. No. 3,306,869 to Lahr et al. discloses floor polishes which are said to possess desirable so-called \"leveling\" properties as well as acceptable \"scuff-resistance\" properties. Such floor polishes comprise a wax, a modified rosin, an emulsifying agent, and a polysiloxane-oxyalkylene block copolymer having a molecular weight below about 25,000 and a polysiloxane content in the range of from about 30 to about 60 percent by weight. U.S. Pat. No. 3,341,338 to Pater also discloses a polish formulation comprising a polysiloxane-oxyalkylene block copolymer and a wax. Neither patent, however, discloses or even suggests \"smear recovery\", or otherwise discloses the specific poly(dimethyl)-co-poly(methylalkyl) siloxane or poly(dimethyl)-co-poly(methyl, oxygen-containing) siloxane copolymers of our present invention.\nFrom the above discussion, it can either be implied that the smear-recovery problems associated with certain prior-art polish compositions are not appreciated or, if such problems have become recognized by those skilled in the art, it can be implied that an acceptable solution to the smear-recovery problem is still generally unknown to those skilled in the art.\nAccordingly there is presently a need for polish compositions which possess rapid smear-recovery properties. One object of our invention, therefore, is to provide polish compositions for furniture and other similar household items, wherein such polish compositions exhibit rapid smear-recovery qualities.\nIt is yet another object of our invention to provide furniture polish compositions which may be applied and removed with a minimum of effort, and which possess not only rapid smear-recovery properties but also acceptable so-called \"depth-of-gloss\" characteristics as well.\nThe foregoing, as well as other objects, features and advantages of our invention will become more readily clear to those skilled in the art upon reference to the following summary and detailed description."}
{"text": "The present invention relates to a swash plate type compressor that has single headed pistons and is used in an air conditioner of a vehicle, and more particularly, to improvement of a radial bearing that supports a drive shaft for reciprocating the pistons and to improvement of a lubricating structure of a shaft sealing assembly.\nAs shown in FIG. 6, the housing of a typical swash plate type compressor includes a front housing member 71, a cylinder block 72 and a rear housing member 73, which are secured to one another. A drive shaft 74 has a first end and a second end. The drive shaft 74 is supported by the housing through a first and second radial bearings 75, 76 such that the first end protrudes from the front housing member 71. A shaft sealing assembly 78 is located about the drive shaft 74 at a position between the first end and the first radial bearing 75. The sealing assembly 78 prevents refrigerant gas from leaking from a crank chamber 77 to the atmosphere.\nMoving parts of a compressor such as bearings are lubricated by misted lubricant contained in refrigerant gas. Therefore, parts where refrigerant gas is stagnant are not effectively lubricated. A compressor that uses carbon dioxide (CO2) for a cooling circuit instead of chlorofluorocarbon has been introduced. When using CO2 as refrigerant, the refrigerant pressure is more than ten times that of a case where chlorofluorocarbon is used as refrigerant, which increases the load acting on bearings and shaft sealing assemblies. Accordingly, lubrication must be improved.\nIn the compressor of Japanese Unexamined Patent Publication No. 11-241681, the shaft sealing assembly 78 is located in an isolated chamber 80, which is forward of the first radial bearing 75. A decompression passage 79 is formed in the drive shaft 74. An outlet 79b of the decompression passage 79 opens to the end face of the second end of the drive shaft 74. A fan 81 is attached to the second end of the drive shaft 74. When the fan 81 rotates integrally with the drive shaft 74, refrigerant in the decompression passage 79 is drawn to the outlet 79b. The refrigerant then flows to the crank chamber 77 through the radial bearing 76.\nThe isolated chamber 80 is connected to the crank chamber 77 through the space in the radial bearing 75 and the space in a thrust bearing 82. The spaces in the radial bearing 75 and the thrust bearing 82 function as oil supplying passages.\nJapanese Unexamined Patent Publication No. 8-165987 discloses a compressor shown in FIG. 7. In this compressor, a second end of the drive shaft 74 faces a chamber 84 that communicates with a suction chamber 83. An axial passage 85 is formed in the drive shaft 74. The inlet 85a of the passage 85 opens to an isolated chamber 80. The outlet 85b of the passage 85 opens to the chamber 84.\nIn the compressor of FIG. 6, the fan 81 attached to the drive shaft 74 draws some of refrigerant gas into the decompression passage 79 through the first radial bearing 75 or through the thrust bearing 82. The drawn refrigerant gas then returns to the crank chamber 77 through the second radial bearing 76. Accordingly, the radial bearings 75, 76 and the shaft sealing assembly 78 are reliably lubricated. However, to flow lubricant through the decompression passage 79, the fan 81 is required, which complicates the structure.\nInstead of a fan, the chamber 84 is located adjacent to the second end of the drive shaft 74 of the compressor shown in FIG. 7, and the passage 85 is formed in the drive shaft 74 to connect the isolated chamber 80 with the chamber 84. Thus, refrigerant flows through the radial bearings 75, 76 or through the thrust bearing 82 in accordance with the pressure difference between the crank chamber 77 and the chamber 84. However, since the inlet 85a is located between the shaft sealing assembly 78 and the thrust bearing, flow of refrigerant is weakened either in the shaft sealing assembly 78 or in the thrust bearing, which results in insufficient lubrication.\nAccordingly, it is an objective of the present invention to provide a swash plate type compressor that includes a simple structure for effectively lubricating radial bearings, which support a drive shaft, and a shaft sealing assembly.\nTo achieve the foregoing and other objectives and in accordance with the purpose of the present invention, a swash plate type compressor is provided. The compressor includes a housing, a drive shaft, first and second radial bearings, a piston, a cam plate, a shaft sealing assembly. A suction chamber, a discharge chamber and a crank chamber are defined in the housing. The housing has at least one cylinder bore. The drive shaft is rotatably supported by the housing and has a first end portion and a second end portion. The first end portion protrudes from the housing. The first and second radial bearings support the first and second end portions of the drive shaft, respectively. The piston is reciprocally accommodated in the cylinder bore. The cam plate is accommodated in the crank chamber and is operably coupled to the piston to convert rotation of the drive shaft into reciprocation of the piston. The shaft sealing assembly seals the space between the drive shaft and the housing and is accommodated in the suction chamber. The suction chamber is closer to the first end portion of the drive shaft than the first radial bearing is. A passage is formed in the drive shaft to connect the suction chamber to the crank chamber. The passage has an inlet and an outlet. The inlet is closer to the second end portion than the second radial bearing is. The outlet is closer to the second end portion than the first radial bearing is.\nOther aspects and advantages of the invention will become apparent from the following description, taken in conjunction with the accompanying drawings, illustrating by way of example the principles of the invention."}
{"text": "The disclosed embodiments generally relate to optical communication links and laser optical ranging, and in particular, to optical communication links and laser optical ranging between extraterrestrial bodies.\nIn space based optical communication and ranging, pointing, or attitude precision is critical as it directly affects the communication link power budget, communication link power fluctuations and the size, weight, and power of the communication or ranging instrumentation. In most examples, the communication link and communication and ranging instrumentation may use optical beams, in particular, lasers. Because of the distances between a space craft and a target, for example, Earth, and the coherence of the optical beam, pointing errors in the micro-radian range may result in missing the target.\nA star tracker may be used to determine the attitude of the space craft with respect to the stars. Because the positions of many stars have been accurately measured over the years, a camera may be used with a database of star locations and patterns to determine and adjust the orientation of the spacecraft. The accuracy of the control systems used for attitude adjustment may limit the pointing accuracy. For example, present reaction wheel attitude control systems may limit pointing accuracy to approximately 50 micro-radians. Attitude dynamics, atmospheric drag, magnetic fields, solar radiation pressure, gravity gradients, vibration, and heating and cooling effects may also affect the ability of a space craft attitude control system to maintain a particular attitude in order to accurately direct a communication or ranging beam toward a target.\nIt would be advantageous to improve provide a system that improves laser beam pointing accuracy to the level of that of the start tracker pointing knowledge so the overall pointing accuracy will not be limited by the accuracy of the control systems used for attitude adjustment."}
{"text": "1. Field of the Invention\nThe present invention relates to an information recording medium drive control apparatus applied to an information processing apparatus provided with a drive that allows the change of information recording mediums.\n2. Description of the Related Art\nIn recent years, information processing apparatuses, such as portable computers, incorporating CD-ROM (Compact Disc Read-Only Memories) drivers have been popularized. CD-ROMs that are capable of driving optical disc drives provided in such information processing apparatuses include music CD-ROMs, photo CD-ROMs, video CD-ROMs, video-game CD-ROMs, car navigation CD-ROMs, and movie CD-ROMs.\nSince these CD-ROMs each have the information recorded in a different format, to read the information from a CD-ROM requires a player compatible with the format of the CD-ROM. To meet this requirement, some commercially available information processing apparatuses of recent date are provided with several types of players so that various types of CD-ROMs can be played back. These players are offered in software form and normally consist of application program modules.\nReferring to a flowchart in FIG. 1, the process of playing back a CD-ROM will be explained.\nFirst, to play back a CD-ROM, the user chooses item \"CD\" on the screen of a display (step A1). This causes a selection menu for prompting the user to choose the type of a CD-ROM to appear on the screen of the display (step A2). The selection menu contains various items including music, photo, video, game, car navigation, and movie, for example. If the user chooses \"PHOTO,\" for example (step A3), the photo player (application program) stored in a ROM or the like is loaded into a main memory (step A4).\nNext, the user inserts a photo CD-ROM into the drive of the electronic apparatus (step A5). The drive drives the CD-ROM automatically and brings it in steady rotation (step A6). Then, the photo player performs the playback process by reading the information from the CD-ROM (step A7).\nAs seen from the processing flow, to play back the desired CD-ROM, the user has to choose a suitable one from the players. Because of this, the user feels that the selection of a player is troublesome. In a case where playback must be performed rapidly, the user may feel that it takes a long time to complete the selection.\nFurthermore, after the selection of the player has finished, the user does not necessarily insert a suitable CD-ROM into the drive. Because the user might insert the wrong CD-ROM by mistake, he or she is always nervous when inserting a CD-ROM, worrying about making a mistake. If the wrong type of CD-ROM is inserted, malfunction may result. When a malfunction has occurred or the playback has not been performed due to the insertion of the wrong type of CD-ROM, the user may waste time until the cause of the malfunction has been corrected.\nRecently, DVDs (digital video discs), whose specification differs from that of conventional CDs, have begun coming onto the market. The types of information recording mediums dealt with by personal computers or the like are expected to increase in number. With many types of players being available, the user will have to choose a suitable one from the many types of players, he or she is more liable to choose the wrong player by mistake, making the selection work more troublesome.\nTo overcome these problems, there have been demands for the techniques for, when an information recording medium is played back, automatically choosing a player compatible with the type of the information recording medium inserted in the drive without the user choosing the player.\nNow, the process of changing the CD-ROM being played back will be described.\nThe photo player is now playing back a photo CD-ROM (step B1). To change the photo CD-ROM to a music CD-ROM, the us e r stops the playback operation of the photo player (step B2) and ejects the photo CD-ROM from the drive (step B3). The user chooses the music player corresponding to music CD-ROMs (step B4) and inserts a music CD-ROM into the drive (step B5). This causes the music player to play back the music CD-ROM (step B6).\nAs seen from the processing flow, to change the photo CD-ROM, the user must stop the playback operation of the photo player. If the user does not stop the player or carries out an operation different from the stop operation by mistake, the following problem will arise. Even after the photo CD-ROM has been changed to a music CD-ROM, the photo player keeps playing pack as it has done. This permits the music player and photo player to read the information from the music CD-ROM at the same time, resulting in malfunction.\nTo overcome this problem, there have been demands for the techniques for, when an information recording medium being played back is changed, automatically stopping t he player corresponding to the information recording medium ejected from the drive without the user stopping the player."}
{"text": "3GPP LTE (3rd Generation Partnership Project Long-term Evolution, hereinafter referred to as “LTE”) uplink uses cyclic shift sequences, which are orthogonal sequences, as pilot signals to reduce interference among sequences. A cyclic shift sequence can be generated by cyclically shifting a pilot sequence by a cyclic shift amount on the time axis. For example, FIG. 1 shows a cyclic shift sequence (m=0) and a cyclic shift sequence (m=1) with pilot sequence length N=12 and cyclic shift amount A=6.\nIn FIG. 1, while the cyclic shift sequence (m=0) is configured in order of a(0) to a(11), the cyclic shift sequence (m=1) is configured, by cyclically shifting the cyclic shift sequence (m=0) by Δ(=6) samples, in order of a(6) to a(11), a(0) to a(5).\nThe cyclic shift amount is determined by a base station apparatus (hereinafter abbreviated to “base station”) and reported from the base station to a terminal station apparatus (hereinafter abbreviated to “terminal”) per scheduling (per subframe). Eight types “0, 2, 3, 4, 6, 8, 9, 10” (3 bits) are defined for reporting the cyclic shift amount. These correspond to a cyclic shift amount of “0, 2, 3, 4, 6, 8, 9, 10”×symbol length/12 (ms).\nSince sequences can be separated with low inter-sequence interference by assigning cyclic shift sequences of different cyclic shift amounts to different terminals, cyclic shift sequences are used for pilot signal transmission in MU-MIMO (Multiple User-Multiple Input Multiple Output). In MU-MIMO, a plurality of terminals transmit data signals at the same time and the same frequency, spatially multiplex the data signals and thereby improve system throughput. At this time, it is also preferable that a plurality of terminals transmit pilot signals at the same time and the same frequency from the standpoint of frequency utilization efficiency. Therefore, cyclic shift sequences, which are orthogonal sequences, for pilot signals and the cyclic shift sequences are transmitted at the same time and the same frequency. The reception side can separate pilot signals using the nature of orthogonal sequences, and can thereby accurately estimate a channel state of each terminal.\nOn the other hand, in LTE-Advanced (hereinafter referred to as “LTE-A”) uplink, studies are being carried out on supporting SU-MIMO (Single User-Multiple Input Multiple Output) to improve throughput, whereby one terminal transmits data signals from a plurality of antenna ports at the same time and the same frequency and spatially multiplexes the data signals using virtual communication channels (hereinafter referred to as “streams”) in the space.\nHere, the “antenna port” refers to a logical antenna (antenna group) made up of one or a plurality of physical antennas. That is, the antenna port does not always refer to one physical antenna, but may also refer to an array antenna made up of a plurality of antennas. For example, the antenna port may be made up of a plurality of physical antennas and defined as a minimum unit whereby a base station or terminal can transmit different pilot signals. Furthermore, the antenna port may also be defined as a minimum unit for multiplying a weight of a precoding vector. Hereinafter, a case will be described as an example where an “antenna port” and a physical antenna have a one-to-one correspondence for simplicity of explanation.\nSU-MIMO requires pilot signals for each stream and studies are being carried out on code-multiplexing pilot signals of each stream using a cyclic shift sequence, which is an orthogonal sequence, for the purpose of reducing inter-sequence interference.\nHere, in an ideal environment in which there is no channel variation, a cyclic shift sequence is an orthogonal sequence and no inter-sequence interference occurs. On the other hand, in a real environment with a channel variation, complete orthogonality is not established and a certain degree of inter-sequence interference occurs. Especially when the number of streams increases and the cyclic shift sequence multiplexing number increases, inter-sequence interference also increases. Therefore, in LTE-A, studies are being carried out on reducing inter-sequence interference using a Walsh sequence as well as cyclic shift sequences adopted in LTE.\nIn multiplexing using Walsh sequences, pilot signals of a first slot (slot #1) and a second slot (slot #2) making up a subframe are multiplied by Walsh sequence w1=[1 1] or Walsh sequence w2=[1 −1] (see FIG. 2). That is, Walsh sequence w1 uses pilot signals similar to those conventional ones in first and second slots and Walsh sequence w2 uses pilot signals similar to those conventional ones in the first slot and uses pilot signals with an inverted phase (180 degree rotation) in the second slot.\nAs a method of reporting a cyclic shift amount, in LTE, the base station reports in three bits using a control information channel (Physical Downlink Control Channel, PDCCH) to be reported to each terminal per scheduling. Furthermore, in LTE-A, studies are being carried out on adding one bit indicating whether a Walsh sequence of each terminal is w1 or w2 using a control information channel (PDCCH), the base station reporting the Walsh sequence to each terminal and each terminal switching between the Walsh sequences.\nFurthermore, in order to reduce inter-sequence interference of cyclic shift sequences between streams in SU-MIMO, Walsh sequence w1 is used for pilot signals of odd-numbered streams and Walsh sequence w2 is used for pilot signals of even-numbered streams (see FIG. 3).\nHere, the “stream number” is a number indicating order in which data is assigned. For example, when data is transmitted with only one stream, suppose a stream transmitted from one antenna port is stream #0 and when data is transmitted with two streams, the stream transmitted from an antenna port different from the above-described port is stream #1. By setting different Walsh sequences depending on whether a stream number is an odd number or even number, it is possible to reduce inter-sequence interference between pilot signals of neighboring streams (see Non-Patent Literature 1). Furthermore, since there is no need for reporting a bit indicating a Walsh sequence, which will be used in the second (stream #1) and subsequent streams, the amount of reporting the cyclic shift amount can be reduced."}
{"text": "The present invention relates to a molding lens with indentation for measuring eccentricity and a method for measuring eccentricity thereof, especially to a lens with at least one indentation that is concentric with optical surface and is arranged on the area outside the optical surface. The indentation is observed by a measuring microscope. Thus by eccentricity test function of the measuring microscope, the eccentricity of the lens is obtained and the eccentric direction is defined so as to modify the lens mould.\nRefer to FIG. 1, a manufacturing process of a molding lens A1 includes following steps: Firstly design front and rear optical surfaces A2 of the lens A1. For example, the optical surfaces A2 are aspheric surfaces having a concave surface and a convex surface. For example, an ultra-precision machine in combination with numerical control (NC) programming is used to define toolpath. That means to set sag value (tool depth) so as to conduct precision machine processing of moulds. As to a curvature lens, the Z-axis is an optical axis of the lens and center of surface is as zero point of X-Y plane. On different positions of the surface of lens, the height difference between a line parallel to the Z-axis and the X-Y plane is the sag value. Moreover, the sag value can also be calculated by equations such as Anamorphic surface, First Type Toric surface or Second Type Toric surface. Then by process of injection molding or die-cast molding, lens are mass produced. After finishing the mould, firstly samples are made to run an eccentricity test for mould modification. Furthermore, each of lens produced needs to go through the test procedure to make sure they are good products.\nRefer to FIG. 2, a conventional eccentricity test of the optical surfaces A2 of the lens A1 is disclosed. A transmission eccentric scale indicator A3 in combination with a rotatable jig A4 is used so as to make a parallel light A5 from a light source on the bottom penetrates through the optical surface of the lens A1 being detected to form an image on a screen A6. The transmission eccentric scale indicator A3 is quite expensive so that the manufacturing cost is increased. The jig A4 clips the lens A1 being detected and carries it to rotate while radius of rotation of a focus A7 on the screen A6 is used as criteria for checking eccentricity. Thus the requirement on eccentricity of the jig A4 itself is quite high and the eccentricity of the jig A4 is no more than 2 μm (10−6 m). Thus the cost for the jig A4 is also quite high due to operation error of the jig that affects the eccentricity test results. Moreover, except the test results, the error of the jig A4 also have effects on quality control of the lens A1. Furthermore, conventional technology can only check eccentric scale of the lens A1 on the screen A6. The eccentricity of the front and the rear optical surfaces A2 can't be distinguished. Therefore, the eccentricity of the lens A1 on the screen A6 can't be used as modification index directly.\nIn other words, it is difficult to decide the eccentricity of the lens A1 derived from errors of which optical surface. Thus the optical surface on the mould can't be modified precisely or efficiently and this leads to trouble on mould modification."}
{"text": "The present invention relates to a method for fabricating integrated circuits and in particular to a method for fabrication of data storage type integrated circuits for protection against errors induced by the bombardment of an ionizing radiation.\nIn the trend towards smaller components in integrated circuits to achieve higher density, an error causing mechanism, not previously considered a factor in larger component circuits, in the form of the passage of alpha particles into the integrated circuit, has become a major problem particularly in those circuits intended for memory storage. The passage of the alpha particles into the integrated circuit generates sufficient hole-electron pairs, by the dissipation of energy, to disrupt the integrity of stored data, as an example, in small memory storage cells requiring only a small number of electrons to be retained for the storage of a binary digit.\nIt has been shown that alpha particles are responsible for charge induced errors. These alpha particles originate from radioactive decay of radioactive impurities present in the materials used to fabricate the integrated circuit. A few parts per million or less of radioactive impurities present in the materials emit a sufficiency of alpha particles, as nuclear decay products, to cause an unacceptably high probability of the electrons generated by an alpha particle decay to change the charge stored in the integrated circuit.\nStated another way, soft errors in storage memories, such as MOS dynamic RAM's, are caused by their exposure to alpha particle radiation. Alpha particles are the natural product of the radioactive decay of uranium and thorium present in most materials used in the fabrication of RAM's and thus radiation produces an unwanted charged particle in the silicon semiconductor material which upsets the wanted charge during the normal device operation.\nAccordingly, it is a prime object of this invention to provide a method of reducing alpha particle induced soft error sensitivity in integrated circuits."}
{"text": "The ultimate goal of a computer (application) program is to control the operation or processing of some form of data. More specifically, the goal of a machine vision system application program is to extract important features from image data from which a description, interpretation, or understanding of the objects in the image can be provided by the machine (computer) for the purpose of classification, reporting and decision making.\nIn a machine vision application, for example, a set of operations or steps are executed in sequence. A typical machine vision program flow diagram is shown in FIG. 1. The steps generally begin with the acquisition of an input image, which is captured by a camera and placed in memory, step 1. The image is optionally preprocessed, step 2, before the image is segmented and features are extracted from it. Segmentation and feature extraction, step 3, of the image into one or more distinct objects can be accomplished, for example by extracting their boundaries. Features can be any data that describes or identifies parts, features (holes, etc.) and objects in the camera scene, such as size, area, length, distance or like parameters.\nFeatures are usually expressed in pixel units that describe the geometric position of a feature in the image relative to some fixed position in the image such as the top left corner of the image (also called an image coordinate system), or in physical units like \"mm\" or inches (also called a camera coordinate system) if the position of the camera in the scene is known or can be measured (usually referred to as calibrating the camera).\nThese features are typically checked against tolerances or classified into various categories, step 4, before a decision, step 5, can be made or pass/fail status and results obtained, step 6. Reports can also be generated, detailing the processing application, step 7. External control inputs, step 8, are used to activate the tolerance checking and to provide nominal tolerance values during execution for a particular feature."}
{"text": "1. Field of the Invention\nThe present invention relates to a transmission power determination method, a communication device and a program.\n2. Description of the Related Art\nDiscussions have been taking place recently regarding secondary usage of a spectrum assigned for primary usage to provide a secondary communication service depending on the use condition of the spectrum. For example, the standard specification for allowing an unused channel contained in a spectrum of the U.S. digital TV broadcast (TV white spaces) to be available for radio communication has been studied in the IEEE802.22 working group (cf. “IEEE802.22 WG on WRANs”, [online], [Searched on Jan. 5, 2009], Internet <URL:http://www.ieee802.org/22/>).\nFurther, according to the report from FCC (Federal Communications Commission) on November 2008, the discussions are directed toward permitting secondary usage of TV white spaces by using a communication device that fulfills a certain condition and has received an authorization. The FCC's report accepts the above-described standard specification of IEEE802.22 which is the pioneering work on the standardization of secondary usage of TV white spaces and further coverts the moves of a new study group in IEEE. Technically, because it is required to perform signal detection at the level of −114[dBm] (SNR is about −19[dB] when NF (Noise Figure) is 11[dB], for example) with use of existing technology, for example, an auxiliary function such as geo-location database access is expected to be necessary (cf. “SECOND REPORT AND ORDER AND MEMORANDUM OPINION AND ORDER”, [online], [Searched on Jul. 10, 2009], Internet <URL:http://hraunfoss.fcc.gov/edocs_public/attachmatch/FCC-08-260A1.pdf>). Further, the FCC is searching for opening a 250 MHz band, which is a part of a 5 GHz band, as a new channel for secondary usage.\nFurthermore, in the EU, there are moves afoot to universally allocate a dedicated control channel called CPC (Cognitive Pilot Channel) for making DSA (Dynamic Spectrum Access) under a long-term strategy. Allocation of CPC is incorporated in the agenda of ITU (International Telecommunication Union)-WP11 in 2011. Technological studies for a secondary usage system that makes DSA are also being progressed in IEEE SCC (Standards Coordinating Committee) 41.\nIn such a background, several research reports have been released recently concerning secondary usage of a spectrum in the case of assuming a broadcasting system, a satellite communication system, a mobile communication system or the like as a primary system. For example, Alan Bok et al., “Cognitive Radio System using IEEE802.11a over UHF TVWS”, Motorola, October 2008 proposes a system architecture in the case of operating a radio system with use of the IEEE802.22 standard on TV white spaces of UHF (Ultra High Frequency). Further, D.Gueny et al.,“Geo-location database technique for incumbent protection in the TV White space”, DySPAN, October 2008 also intends use of TV white spaces and proposes a form that utilizes positional information of a service area of a primary system as external information.\nOn the occasion of secondary usage of a spectrum, it is generally necessary for a system on the part of secondary usage (secondary system) to carry out the operation that does not degrade the communication quality of a primary system. Therefore, when transmitting a radio signal in the secondary system, it is desirable to control its transmission power so as to avoid interference on a node of the primary system.\nRegarding such control of a transmission power, in the case of secondary usage of TV white spaces as proposed by Alan Bok et al. or D.Gueny et al., it can be confirmed beforehand that a channel for secondary usage is not used at all, and it is thus possible to determine in many cases that a transmission power at the maximum level can be used. On the other hand, H.Fujii and H.Yoshino (NTT docomo), “Spectrum sharing by adaptive transmit power control for low priority system and its achievable capacity”, CrownCom, May 2008 proposes a technique that protects a node of a high-priority system by adaptively controlling a transmission power in a low-priority system.\nFurther, Inage et al., “Spectrum Sharing Based on Capacity Conservation Ratio of Primary User”, IEICE Technical Report SR2009, May 2009 proposes a technique that, when a system such as a mobile communication system in which the receiving environment of a terminal varies depending on location due to fading or the like is the primary system, adopts the ratio of capacity (capacity conservation ratio) between before and after secondary usage in the primary system as a protection criterion and makes transmission power control for satisfying the capacity conservation ratio."}
{"text": "To establish new turf it is usually necessary to apply fertiliser and also to water regularly. A greater amount of fertiliser and water is required to establish new turf than to maintain existing turf. The type of fertiliser used and the manner of irrigation will affect the rate of turf growth, the structure of the soil and the root structure of the turf as it grows.\nAppropriate fertilisers for turf advantageously comprise a source of nitrogen, phosphorous and potassium. These nutrients may be found in an inorganic or synthetic fertiliser, or in an organic fertiliser.\nWhile inorganic or synthetic fertilisers are often utilised, over application of such fertilisers may be toxic to turf, due to their water solubility. This may also lead to damage to the environment, through leaching of the fertiliser into waterways.\nIn contrast, over application of organic fertilisers is not usually detrimental to turf, as the water soluble content of such fertilisers is normally quite low. Organic fertilisers may also biodegrade over an extended period of time, improve moisture retention and promote microbial and earthworm activity.\nThe water applied to turf may be lost through evaporation, transpiration (emission of water vapour from the leaves of the turf), infiltration and percolation of water through the soil, and water runoff. Due to variations in climate and soil structure, the exact amount of water that must be applied to establish or maintain turf will vary. If the consumer applies too little water, the turf may fail. However, excessive watering will lead to waste.\nNew turf should be watered each day for adequate establishment. The amount of water required may be substantial, and in some cases 4 mm of water may be required per day in the first five weeks after laying to establish new turf. Australian turf production per annum is 6500 hectares (65 million m2), and therefore to bed all of this turf would require on the order of 9.1 mega liters in the first five weeks after laying.\nTo decrease the amount of water required, products such as water absorbent polymers (water crystals) or soil wetting agents are often used. Soil wetting agents help soils to absorb water by decreasing their hydrophobicity. In contrast, water absorbent polymers absorb and retain water for the roots of turf to use. When these polymers are present in soil, it is not necessary to water as frequently as water loss is diminished.\nSubstantially more fertiliser is generally added to soil to establish new turf than other products, such as water absorbent polymers or soil wetting agents. Accordingly, it can be time consuming and difficult to apply the correct proportions of each product to the soil.\nA traditional method of applying multiple products simultaneously would be to simply mix the products together to form a “salt and pepper” mixture. However, during the course of production, distribution and/or application, settling may occur, which in turn limits the ability of each product to be applied at the correct rates due to phase separation based on particle size."}
{"text": "A computer keyboard and display are common devices for providing computer input and output, respectively. A keyboard is language-specific such that the keys on the keyboard can be pressed to directly input only those characters in the keyboard's language that are assigned to those keys. For inputting other characters, a user has to press a combination of keys on the keyboard.\nMany languages have alphabets that are too large to accommodate on a keyboard. Many languages have modifiers, which when applied to a character in the language's alphabet produce additional characters in the language's alphabet. Furthermore, the alphabets of many languages do not use characters to form words in the manner of the English language, but have a collection of characters that represent words. Thus, providing computer input in many languages is not as simple as pressing the letter-keys on the keyboard but an indirect process of pressing a combination of keys to generate characters not available as keys on the keyboard.\nA display is also language-specific. In order to display characters in certain languages, certain scripts, fonts, and rendering directions can be selected for proper rendering of the characters. Generally, for rendering characters of a given language, specific fonts are installed and presented on the computer display.\nText in digital content is encoded in a variety of ways for capturing from an input method, storage in a data storage system, and rendering on a user interface in a desirable manner. For example, Unicode is a common encoding standard for encoding and handling of text characters from a set of language scripts in digital textual content.\nPresently, Unicode encodes or codifies over one hundred thousand characters from over one hundred languages. Generally, in Unicode and other character encoding standards, each character represented in the standard is assigned an encoded value called a code point. For example, in Unicode standard UTF-8 uses one byte—eight bits—to encode the code points of the characters represented in UTF-8. UTF-16 similarly uses two bytes for encoding the code points of the characters represented therein.\nA code point in an encoding standard is unique to a specific character in a specific language represented in that encoding standard. For example, in Unicode, a code point comprises an alphanumeric string that can be generated on commonly used keyboard configurations, such as an English language QWERTY keyboard.\nAs an example, to provide a Unicode code point, the user or an application generally supplies an indication that the alphanumeric string following the indication is a Unicode code point and should be translated using a Unicode table to generate a character. For example, in some implementations, to provide a Unicode code point using a QWERTY keyboard, a user can press the ALT key, keep the ALT key depressed while entering the code point, and release the ALT key when the code point entry is complete.\nAn application called an input method application, or simply “input method” or “IMA”, intercepts the Unicode code point that the user enters. The IMA looks up a Unicode table to find the character that matches the code point that the user entered. The IMA supplies the character to a target application for which the user is supplying the input.\nA character variant is an alternate representation of a character within the given encoding standard. There are two types of character variants—Variants with Different Code Points (VDCP) and Glyphs with Same Code Point (GSCP). A variant is a VDCP type variant when the original character and the variant have different code points, with the same or different language tags, within the encoding standard. A variant is a GSCP type variant when the glyph—a manner of scripting or otherwise visually representing the character—of the original character and the variant has the same code point but different language tags in the encoding standard."}
{"text": "In a communication network system such as the UMTS terrestrial radio access network (UTRAN), there are two potential bottlenecks, namely, the air interface and the transport network (transport link) connecting the radio network controller (RNC) and Node B. The transport link between the RNC and Node B is a potential bottleneck when its capacity is smaller then the available maximal capacity of the Uu interface. For example, a typical scenario is that the Node B is connected to the RNC through an E1 link with a capacity of approximately 2 Mbps, and in this case the available Uu capacity for the high speed downlink packet access (HSDPA) may be significantly larger that 2 Mbps. This means that a single user equipment (UE) with good radio conditions can overload the transport network (TN).\nThe fair sharing of Uu resources is the task of the Uu scheduler, but the Uu scheduler can not cope with the TN bottleneck, i.e. the transport link bottleneck. In order to deal with the TN bottleneck a flow-control (FC) mechanism has been introduced. FIG. 2 shows the location of the FC in the protocol stack. The goal of the FC is to efficiently use the TN in a fair manner.\nLack of FC causes serious performance degradation when the transport network is the bottleneck. In this case the TN buffer is typically full, causing high TN delay and loss ratio. This causes exhaustive radio link control (RLC) retransmissions which results in a much lower throughput. In addition to this, an RLC reset and a consequent transmission control protocol (TCP) timeout may also occur.\nThe flow control operates per-flow basis, i.e. each HSDPA flow has its own (i) congestion detection, (ii) bitrate calculation and (iii) shaper part. The main tasks of these three parts are the following:\n1) Congestion Detection Part in the Node B\nBased on the arrived packets from the RNC there is an attempt to determine the congestion level of the transport network. If TN congestion is detected, the bitrate calculation part is informed. A gap in sequence numbers of arriving packets is interpreted as “hard” congestion, because with a very high probability this event is due to packet loss in the TN caused by serious congestion. In addition to this, the variation of the one-way packet delay between RNC and Node B is also measured, i.e. a given fraction of packets have a time-stamp. If this delay starts to increase, probably due to queue build up in the TN, then it is interpreted as “soft” congestion, but if this delay build up is getting too large, e.g. larger than 60 ms, it is interpreted as “hard” congestion. The bitrate calculation part will react on hard and soft congestions in different ways.\n2) Bitrate Calculation Part in the Node B\nThis part of the flow control calculates the current maximum bitrate of the flow. This bitrate is allowed by the transport network for that flow. The applied algorithm is conformed with the additive increase multiplicative decrease (AIMD) property that guarantees convergence to fairness; all flows converge to an equal share of resources in steady state, where no flows join or leave. The FC maintains an internal variable for the maximum bitrate of the flow. This bitrate is increased linearly if there is no TN congestion, i.e. no reported congestion from congestion detection part. If congestion is reported, the bitrate is reduced by 50% in case of hard congestion and reduced by 10% in case of soft congestion. When a new flow arrives, in this way a new FC entity is created, a slow-start like mechanism is used to find out the proper starting bitrate of the flow. After the first congestion, the FC behaves the above described AIMD manner. If the calculated bitrate of the flow changes significantly, then the shaper is informed about the new bitrate through a control frame called a capacity allocation (CA). To avoid too high processing load, this part of the FC is executed periodically with a 100-ms period, i.e. the bitrate calculation part is executed every 100 ms.\n3) Flow Shaper in the RNC\nThe task of the shaper is to shape the flow according to the signalled maximum flow bitrate. This bitrate is coming from the latest received CA control frame.\nThe current HSDPA flow control solution provides fairness only in long term due to the convergence of AIMD property in case of a TN bottleneck. Fairness is provided only among flows sharing the same TN bottleneck. The initial shaping rate calculated by the FC has significant effect of the fairness and the time of the convergence. If a new flow arrives into the cell, e.g. due to handover, the existing algorithm operates as follows:\nFirstly, it estimates maximum achievable bandwidth on the transport network (maxHsRate) and the maximum achievable peak rate for HS in the given cell (maxUuRate).\nSecondly, it counts the number of active flows belonging to the Node B (nPqsRbs) and to the cell (nPqsCell). Noticeable is that the new flow also is counted in these counters.\nFinally, it calculates the average bitrate on the transport network and in the cell by dividing the bitrates by the number of ongoing flows respectively. Then it chooses the minimum of them as initial shaping rate of the new flow. Additionally, there is an upper limit (hsSsStartPointMax) for the initial shaping rate to avoid too high initial rates.\nIn another words, this calculation estimates the theoretical fair bandwidth share from the estimated maximum available bandwidth assuming the system is fully utilized and the flows share it equally. Then it sets this as the starting point of the new flow. For instance, if the maximum achievable bandwidth on the transport network is 2 Mbps, the maximum achievable peak rate in the cell is 3.6 Mbps, on the transport there are 5 parallel flows, but on the cell there are only 3 flows and the hsSsStartPointMax is 500 kbps then, the initial shaping rate of the new flow is:CAinitial=min(2 Mbps/6; 3.6 Mbps/4; 500 kbps)=333 kbps  (1)\nNoticeable is that this calculation contradicts the per-flow manner, since the number of flows is aggregated information about the system. In case of steady state, when the HSDPA flow control entities have enough time to find the fair share of the flows, and there is unused capacity neither in the transport network nor in the cell, the above introduced methodology provides good estimation of the fair share, so the convergence will be very fast.\nHowever, the above described solution requires an estimation of the maximum achievable transport network bitrate. One option for this estimate is a parameter configured based on knowledge about the transport network architecture.\nFurther, the initial bitrate calculation relies on rough estimations and, the initial bitrate calculation does not take fairness into account.\nIf the estimation of the maximum achievable bandwidth on transport network (maxHsRate) or the maximum achievable peak cell rate (maxUuRate) is not accurate, they typically overestimates the real ones, the initial shaping rate will be quite far from the optimal and also in terms of fairness. Typically the actual available bitrate for HSDPA is smaller than the maximum, due to the bitrate used by higher priority traffic."}
{"text": "Spectroscopy is a common technique for measuring the concentration of organic and some inorganic constituents of a solution. The theoretical basis of this technique is the Beer-Lambert law, which states that the concentration ci of an absorbent in solution can be determined by the intensity of light transmitted through the solution, knowing the pathlength dλ, the intensity of the incident light I0,λ, and the extinction coefficient εi,λ, at a particular wavelength λ. In generalized form, the Beer-Lambert law is expressed as:\n                              I          λ                =                              I                          0              ,              λ                                ⁢                      ⅇ                                          -                                  d                  λ                                            ·                              μ                                  a                  ,                  λ                                                                                        (        1        )                                          μ                      a            ,            λ                          =                              ∑                          i              =              1                        n                    ⁢                                    ɛ                              i                ,                λ                                      ·                          c              i                                                          (        2        )            where μα,λ is the bulk absorption coefficient and represents the probability of absorption per unit length. The minimum number of discrete wavelengths that are required to solve EQS. 1-2 are the number of significant absorbers that are present in the solution.\nA practical application of this technique is pulse oximetry, which utilizes a noninvasive sensor to measure oxygen saturation (SpO2) and pulse rate. In general, the sensor has light emitting diodes (LEDs) that transmit optical radiation of red and infrared wavelengths into a tissue site and a detector that responds to the intensity of the optical radiation after absorption (e.g., by transmission or transreflectance) by pulsatile arterial blood flowing within the tissue site. Based on this response, a processor determines measurements for SpO2, pulse rate, and can output representative plethysmographic waveforms. Thus, “pulse oximetry” as used herein encompasses its broad ordinary meaning known to one of skill in the art, which includes at least those noninvasive procedures for measuring parameters of circulating blood through spectroscopy. Moreover, “plethysmograph” as used herein (commonly referred to as “photoplethysmograph”), encompasses its broad ordinary meaning known to one of skill in the art, which includes at least data representative of a change in the absorption of particular wavelengths of light as a function of the changes in body tissue resulting from pulsing blood. Pulse oximeters capable of reading through motion induced noise are available from Masimo Corporation (“Masimo”) of Irvine, Calif. Moreover, portable and other oximeters capable of reading through motion induced noise are disclosed in at least U.S. Pat. Nos. 6,770,028, 6,658,276, 6,157,850, 6,002,952 5,769,785, and 5,758,644, which are owned by Masimo and are incorporated by reference herein. Such reading through motion oximeters have gained rapid acceptance in a wide variety of medical applications, including surgical wards, intensive care and neonatal units, general wards, home care, physical training, and virtually all types of monitoring scenarios."}
{"text": "The invention relates to a coupler for a pipe or hose section.\nU.S. Pat. No. 4,407,532 discloses a coupler with a hose section having a coupling body with a coupling piece integrally formed therewith and a pressing sleeve which can be pre-assembled to the coupling body. When using this coupler, a connectable hose end is introduced in an annular gap between the coupling piece and the pressing sleeve and connected by concentrically pressing the pressing sleeve together with the coupler.\nThe desired pre-assembly between the coupling body and the pressing sleeve is achieved with known couplers by inserting the pressing sleeve with its front end over the area of the coupling body adjacent to the coupling piece and pressing a stop projection projecting inwardly from the pressing sleeve through a plastic deformation in an annular stop groove in the coupling body. Since the inside diameter of the stop projection of the pressing sleeve is larger than the outside diameter of the annular projection, in order to be able to insert the same over the annular projection formed between the stop groove and the coupling piece, prior to achieving the pre-assembly, substantial pressure must be applied to the plastic deformation to achieve the pre-assembly and a special tool is required for this. The rejection rate with couplers is relatively high in the production of such pre-assembly. The above-mentioned pre-assembly between the coupling body and the pressing sleeve is not suitable for the use of coupling bodies made of plastic since they would not withstand the pressures exerted upon them when producing the pre-assembly. The use of coupling bodies of plastic is, however, advantageous with regard to use for coupling for pipes or hoses carrying water.\nThe invention therefore relates to a generic coupling in which the coupling body and the pressing sleeve can be pre-assembled in a simple manner that is favorable in terms of production costs.\nThrough the use of a locking or interlocking connection, the pre-assembly of the coupling body and pressing sleeve can be achieved rapidly and simply without tools and without risk of damage.\nWhen selecting a locking connection, it is preferable that the front end be designed in an elastically deformable manner with the pressing sleeve where the front end of the pressing sleeve preferably locks with the coupling body upon spring-back.\nThe inventive coupler is advantageously suited for the use of plastic, particularly polymers such as polysulfone or polyphenylene sulfone, as the material of the coupling body where the pressing sleeve can consist of stainless steel.\nWith such a coupler, pipe or hose sections can be connected with as many additional connections as desirable such as additional pipe or hose sections. Other arrangements such as Ts, reducers, elbows, transitions to sections with internal or external threads or similar types are possible with the inventive coupler. Further advantages and details of the invention are shown in the accompanying drawings. In the Figures:"}
{"text": "The invention generally relates to systems and methods for guiding or locating diagnostic or therapeutic electrode elements in the interior regions of the body. More particularly, the invention relates to guiding or locating diagnostic or therapeutic electrode elements inside the heart for treatment of cardiac conditions.\nPhysicians make use of catheters today in medical procedures to gain access into interior regions of the body f or diagnostic and therapeutic purposes. It is important for the physician to be able to precisely position the catheter within the body to gain contact with a desired tissue location.\nThe need for precise control over the catheter is especially critical during procedures that ablate myocardial tissue from within the heart. These procedures, called ablation therapy, are used to treat cardiac rhythm disturbances.\nDuring these procedures, a physician steers a catheter through a main vein or artery into the interior region of the heart that is to be treated. The physician then further manipulates a steering mechanism to place the electrode carried on the distal tip of the catheter into direct contact with the endocardial tissue. The physician directs energy from the electrode through myocardial tissue either to an indifferent electrode (in a uni-polar electrode arrangement) or to an adjacent electrode (in a bipolar electrode arrangement) to ablate the tissue.\nBefore ablating heart tissue, physicians often examine the propagation of electrical impulses in heart tissue to locate aberrant conductive pathways and to identify the arrhythmia foci, which are ablated. The techniques used to analyze these pathways and locate foci are commonly called xe2x80x9cmapping.xe2x80x9d\nConventional cardiac tissue mapping techniques use multiple-electrodes positioned in contact with heart tissue to obtain multiple electrograms. These conventional mapping techniques require invasive open-heart surgical techniques to position the electrodes on the interior or exterior surfaces of the heart.\nAn alternative technique of introducing multiple-electrode arrays into the. heart through venous or arterial accesses to map myocardial tissue is known. Compared to conventional, open heart mapping techniques, endocardial mapping techniques, being comparatively noninvasive, hold great promise. Multiple electrogram signals obtained from within the heart can be externally processed to detect local electrical events and identify likely foci.\nOnce mapping identifies the foci, an ablation electrode is steered into position in contact with a focus site. At least, in theory, this is the goal sought to be achieved. In actuality, though, the task of remotely guiding an ablation element within the blood pool of a beating heart to a particular focus site can be, at best, problematic.\nThere is the need to provide simple, yet reliable, ways of guiding electrode elements within the, heart, or within other interior parts of the body, to precise locations targeted for diagnosis or treatments.\nThis invention has as its principal objective the realization of safe and efficacious systems and methods for remotely locating electrode elements at precise locations within the body.\nOne aspect of the invention provides a system and related method for guiding a movable electrode within an array of multiple-electrodes located within the body. The system and method couple an emitting electrode to an electrical energy-generating element. The emitting electrode comprises either the movable electrode or at least one of the electrodes in the array. The generating element conditions the emitting electrode to emit electrical energy while the movable electrode is located within the array.\nAccording to this aspect of the invention, the system and method couple a sensing electrode to a sensing element. The sensing electrode comprises either the movable electrode or at least one of the electrodes in the array. The sensing element conditions the sensing electrode to sense electrical energy emitted by the emitting electrode.\nFurther in accordance with this aspect of the invention, the system and method couple a processing element to the sensing element. The processing element analyzes sensed electrical energy and generates, based upon its analysis, an output that locates the movable electrode within the array.\nAnother aspect of the invention provides a system and method for ablating tissue within the heart. The system and method are usable in conjunction with an array of multiple-electrodes located within the heart to locate foci appropriate for ablation and an ablation electrode that is movable within the array for applying ablating energy to the foci.\nAccording to this aspect of the invention, the system and method condition an emitting electrode to emit electrical energy while the ablation electrode is present within the array. The emitting electrode comprises either the ablation electrode or at least one of the electrodes in the array. While the ablation electrode is present within the array, the system and method also sense the emitted electrical energy with a sensing electrode. The sensing electrode comprises either the ablation electrode or at least one of the electrodes in the array. The system and method process the sensed electrical energy to generate an output locating the ablation electrode relative to the multiple-electrodes on the array.\nIn a preferred embodiment, the system and method continuously generate the location indicating output while the physician moves the ablation electrode within the array. In this way, the system and method aid the physician in guiding the ablation electrode to the targeted ablation site.\nIn a preferred embodiment, the systems and methods that incorporate either aspect of the invention generate an electric field within the array, while sensing electrode electric potentials in the electric field. In this embodiment, the processing element generates the output by analyzing spatial variations in the electrical potentials within the array. The variations can comprise variations in phase, variations in amplitude, or both. Alternatively, the processing element generates the output by analyzing spatial variations in impedances between the emitting and sensing electrodes.\nIn a preferred embodiment, the systems and methods that incorporate either aspect of the invention inject electrical energy into body tissue, while sensing tissue response to the injected electrical energy. In. this embodiment, the processing element generates the output by analyzing differences in the sensed tissue response. In one implementation, the processing element analyzes time differences in the sensed tissue response. In another implementation, the sensing element senses the depolarization of myocardial tissue, and the processing element analyses time differences in sensed tissue depolarization.\nYet another aspect of the invention provides a system and method for ablating tissue within the heart. The system and method locate an array of multiple-electrodes in contact with tissue within the heart to sense electrical activity in heart tissue to locate foci appropriate for ablation. The system and method also locates a movable ablation electrode within the array. The system and method condition an emitting electrode comprising either the ablation electrode or at least one of the electrodes in the array to emit ultrasonic energy while the ablation electrode is located within the array. While the ablation electrode is located within the array, the system and method also sense the emitted ultrasonic energy with a sensing electrode. The sensing electrode comprises the ablation electrode, if the ablation electrode is not the emitting electrode, or, otherwise, at least one of the electrodes in the array. The system and method process the sensed ultrasonic energy to generate an output locating the ablation electrode relative to the multiple-electrodes in the array.\nIn a preferred embodiment of this aspect of the invention, the system and method move the ablation electrode within the array while repeating the emitting, sensing, and processing steps just described. The result is an output that continuously locates the ablation electrode as it moves within the array.\nOther features and advantages of the inventions are set forth in the following Description and Drawings, as well as in the appended Claims."}
{"text": "1. Field of the Disclosure\nThis disclosure relates to tightening systems for use in fitting a wearable article, such as a helmet.\n2. Description of the Related Art\nHelmets and other wearable articles are commonly used to provide protection to the head or other body parts of a wearer, such as during sporting activities and other activities. Some helmets comprise a hard shell of plastic or Kevlar® or the like, and various pads, straps and bladders to position the helmet on a wearer's head. Heads come in a wide variety of shapes and sizes. Helmet shells, on the other hand, come in a very limited number of sizes. If a helmet does not fit properly to the wearer's head, it can cause discomfort and may not provide sufficient protection in some cases. For example, if a helmet is worn that is too large for the wearer's head, the helmet can shift positions during use and may even fall off. Helmets can be made to fit a variety of head sizes and shapes, but existing helmets suffer from various drawbacks. For example, some existing helmets do not provide sufficient adjustability to comfortably fit to a wide variety of head shapes and sizes. Some existing helmets apply pressure unevenly across the head of the wearer, which can cause discomfort."}
{"text": "A flash or block erase memory (flash memory), such as, Electrically Erasable Programmable Read-Only Memory (Flash EEPROM), includes an array of cells which can be independently programmed and read. The size of each cell and thereby the memory as a whole are made smaller by eliminating the independent nature of each of the cells. As such, all of the cells are erased together as a block.\nA memory of this type includes individual Metal-Oxide Semiconductor (MOS) memory cells that are field effect transistors (FETs). Each FET, or flash memory cell includes a source, drain, floating gate and control gate to which various voltages are applied to program the cell with a binary 1 or 0, or erase all of the cells as a block. The flash memory cell provides for nonvolatile data storage.\nA typical configuration of a dual bit flash memory cell consists of an oxide-nitride-oxide (ONO) layer, within which the storage element is contained. The ONO layer is sandwiched between a control gate and a crystalline silicon semiconductor substrate. The substrate includes selectable source/drain regions.\nIn a dual bit flash memory cell, the flash memory cell stores data by holding charge in the ONO layer. The charge storage element within the ONO layer allows electrons to be stored on either side of the flash memory cell. That is, charge is stored within the ONO layer on either side of the mirror bit memory cell. As a result, the basic memory cell behaves as two independent conventional memory cells. In a typical dual bit flash memory cell, a program operation is done by injecting hot electrons into the ONO layer and an erase operation is done by injecting hot holes into the ONO layer.\nA typical configuration of a floating gate flash memory cell consists of a thin, high-quality tunnel oxide layer sandwiched between a conducting polysilicon floating gate and a crystalline silicon semiconductor substrate. The tunnel oxide layer is typically composed of silicon oxide (SiO). The substrate includes a source region and a drain region that can be separated by an underlying channel region. A control gate is provided adjacent to the floating gate, and is separated by an interpoly dielectric. Typically, the interpoly dielectric can be composed of an oxide-nitride-oxide (ONO) structure.\nThe flash memory cell stores data by holding charge within the floating gate. In a write operation, charge can be placed on the floating gate through hot electron injection, or Fowler-Nordheim (F-N) tunneling. In addition, F-N tunneling can be typically used for erasing the flash memory cell through the removal of charge on the floating gate.\nA typical configuration of an array of dual bit or floating gate flash memory cells includes rows and columns of flash memory cells. The array is supported by word lines and bit lines, wherein the word lines are coupled to gates of flash memory cells, and the bit lines are coupled to source and/or drains.\nAs the flash memory cells within an array are scaled to higher values, the size of structures within the array are decreased. For example, the bit lines become narrower and narrower within the array of memory cells. However, as the bit lines become narrower, misalignment issues between bit line contacts and the bit lines become more pronounced. As a result, the alignment margin when forming the contact to the bit lines becomes very small, thereby limiting the scaling increase of the array of memory cells.\nIn addition, in the conventional art, the bit lines can not be silicided, such that, a cobalt silicide (CoSi) layer cannot be formed on the bit lines. The CoSi layer provides for better conductivity, for example between the contact and the bit line. That is, the CoSi layer lowers the bit line contact resistance. However, in conventional manufacturing techniques forming the CoSi layer would electrically couple all the bit lines together in the array of memory cells. In that case, individual cells could not be isolated for reading or programming, since any bit line is coupled to all of the memory cells in the array.\nAs a remedy for the CoSi short, ONO layers are not removed between the word lines and in the bit lines contact regions. As such, when a CoSi layer is deposited for the gates of the memory cells and for the source/drain regions of the periphery, the ONO acts to block CoSi formation between the bit lines of the array of memory cells. However, keeping the ONO layer between the word lines also can induce a leakage current. As a result, a leakage current would cause an entire column of memory cells to malfunction. In particular, with the leakage current, the total current read from a column of memory cells will include the current from the programmed memory cell being read and the leakage current. This may result in the programmed cell appearing to be erased."}
{"text": "The present invention relates to a liquid crystal panel type projection display More particularly, the present invention is directed to a compact, high definition, rear projection type display liquid crystal panel display with an improved lens or liquid crystal panel control system.\nThe construction of a conventional liquid crystal panel type rear projection display is shown in FIG. 2 wherein the numeral 1 denotes a light source such as a halogen lamp or a xenon lamp, numeral 2 denotes a condenser mirror, numeral 3 denotes a dichroic mirror having the property of reflecting red color and transmitting green and blue colors (red, green and blue colors will hereinafter be referred to simply as R, G and B, respectively), and numeral 4 denotes a dichroic mirror for reflecting B and transmitting R and G. Numerals 5, 6, 7 and 8 denote mirrors, and numerals 9, 10 and 11 denote liquid crystal panels for R, G and B, respectively. The liquid crystal panels may be so-called TFT (Thin Film Transistor) type nematic liquid crystal panels. Numerals 12, 1 3 and 20 denote a dichroic prism, a projection lens and a screen, respectively.\nIn the construction shown in FIG. 2, color images can be projected onto the screen 20 by inputting image signals to the liquid crystal panels 9, 10 and 11. Such prior art construction is described in '86 Television Society Report IPD 109-5, pp. 23-28. Further, Japanese Patent Application Laid-Open No. 19834/87 describes the use of a Fresnel lens for a liquid crystal panel for controlling light input or output therefrom, but is not concerned with viewing angle of a projection system.\nThe above-noted prior art is advantageous in that the use of only one projection lens suffices On the other hand, however, it has the drawback that it is structurally difficult to set the value of an image projection field angle, .alpha., shown in FIG. 2 at 25.degree. or more. The value .alpha. in the prior art is about 23.degree..\nFor a front projector, a field angle of 23.degree. or so is appropriate, but for constituting a compact, rear projection type display it is considered necessary to use an optical system having a field angle of about 25.degree. or more, preferably 35.degree. or more. There are two reasons why it is difficult to obtain a field angle of 25.degree. or more in the conventional one-lens, three-panel type construction shown in FIG. 2. The first reason is based upon a physical restriction with regard to the presence of the dichroic prism 12. The second reason is that since the electrooptical conversion characteristic per se of the liquid crystal panels depend on the direction of exit light from the panels, there occurs unevenness in luminance and in color, causing the image quality to be impaired, at diagonal corners on the screen in an optical system having a large field angle of 25.degree. or more."}
{"text": "1. Field of the Invention\nThe invention relates to a system on chip (SOC), and more particularly, to isolation circuits in a level shifter.\n2. Description of the Related Art\nFIG. 1 shows a conventional system on chip (SOC) utilizing a level shifter 100. Signals passed between the first voltage domain 110 and the second voltage domain 120 are converted by the level shifter 100 to adapt to the corresponding voltage domain. There are various known techniques for implementing level shifter 100, thus, a detailed description thereof is omitted. In an 500, however, the first voltage domain 110 or second voltage domain 120 may be selectively or occasionally powered down, switching the corresponding terminal on level shifter 100 to a floating state. For example, when the first terminal V1 is floating, the level shifter 100 may output an indeterminable signal to the second terminal V2. If the second voltage domain 120 receives the indeterminable signal, unexpected application errors may occur."}
{"text": "This invention relates to antennas and more particularly to a low physical profile antenna that can be arrayed to provide near isotropic radiation patterns.\nIn the past, numerous attempts have been made using stripline antennas to provide an antenna having ruggedness, low physical profile, simplicity, low cost, and conformal arraying capability. However, problems in reproducibility and prohibitive expense made the use of such antennas undesirable. Older type antennas could not be flush mounted on a missile or airfoil surface. Slot type antennas required more cavity space, and standard dipole or monopole antennas could not be flush mounted."}
{"text": "The invention relates to improved mechanisms for the communication of network messages between two network nodes on a physical connection being established between those nodes at a physical layer switch.\nIt is common practice for the release of news bulletins relating to certain events to be tightly controlled such that the bulletins are released not before a scheduled time. This is particularly well known in finance where there can be significant value in being the first to have access to information that could affect the markets. For example, news bulletins carrying the details of governmental budgets or interest rate changes by central banks are often embargoed until a predetermined time, at which point the news agencies that have been permitted to report on the event simultaneously release their bulletins into the public domain.\nNews bulletins are typically embargoed until a predetermined time by providing press reporters with access to the news information only within a “lock-up room” that is isolated from the outside world, with no communications being permitted from the room. Within the lock-up room, the news reporters are free to draft bulletins reporting the news event on computers provided for that purpose. However, those computers are physically isolated from public communication networks by an “air gap”. Network messages carrying the news bulletins are therefore queued for delivery at the transmit queues of the computers until, at the predetermined time, a switch is thrown and a physical connection to the news distribution network is established. Such press lock-ups are used by the Australian and Canadian governments so as to provide for a scheduled release of Federal Budget information (see http:(slash)(slash)www2b.abc.net.au/guestbookcentral/entry.asp?GuestbookID=389&EntryID=7 55777 and http:(slash)(slash)www.cbc.ca/news/background/budget2006/blog.html), as well as by the US Department of Labor.\nOn a connection being made from the computer to the news distribution network, the physical and logical links appropriate to the communication protocols in use at the computer must be established so as to permit the transmission of the queued messages onto the network. For example, when a computer in the lock-up room is reconnected to the outside world, the computer would typically establish a connection with a server located outside of the lock-up room and configured to provide a gateway onto the respective news provider's network. If the connection were an Ethernet connection, then a physical layer link must first be established between the computer and server, over which a logical data link can subsequently be established to the intended endpoint receiver of the messages from that computer. For 100BASE-T Ethernet, the time required to establish such a data link layer connection can be 100 ms or more.\nThe advent of high speed trading has meant that significant profits can be made by traders who are able to exploit microsecond advantages in the receipt of financial information. Delays of tens of milliseconds therefore represent a significant length of time. Furthermore, the physical switch by which the computers of a lock-up room are isolated from the public networks will not close all of its ports simultaneously. There is typically a random distribution in its port closure timings with millisecond order standard deviation such that the time between a given pair of ports closing can be significant. This inadvertently causes the news bulletins from the lock-up computers allocated to some news providers to be released prior to the bulletins of other providers.\nThere is therefore a need for an improved mechanism for the scheduled release of embargoed news bulletins, particularly those bulletins carrying financial news."}
{"text": "1. Field of the Invention\nThe present invention relates to the design of electrical/electronic circuits, and more specifically to a method and apparatus for attaining a bandwidth limited sampling circuit of high linearity.\n2. Related Art\nA sampling circuit generally refers to a component which samples a signal level of an input signal (e.g., analog signal) at a particular instant of time and provides an output signal having the signal level to other components for a long duration. In general, a sampling circuit samples an input analog signal at a time point specified by a clock signal, and provides an output signal with the sampled strength to other components for further processing for a long time.\nAn input signal generally contains a source signal and a noise signal/component. It is generally desirable that the sampled output generated by a sampling circuit contain only the source signal component, which is generally referred to as providing a high signal to noise ratio (SNR).\nA prior approach may remove the noise component by eliminating high frequency components from an input signal as the noise component is generally present at high frequencies. One problem with such an approach is that some of the components performing such removal may have non_linear characteristics, which results in a non_linear response in terms of the overall sampling operation. The non_linearity is undesirable in several environments. Accordingly, what is needed is a bandwidth limited sampling circuit of high linearity providing high SNR.\nIn the drawings, like reference numbers generally indicate identical, functionally similar, and/or structurally similar elements. The drawing in which an element first appears is indicated by the leftmost digit(s) in the corresponding reference number."}
{"text": "The present invention relates to the preparation of novel diamine antiviral agents containing an aromatic nucleus and their use as prophylactic and therapeutic agents in vertebrates. It is believed that these compounds function as endogeneous interferon inducers though the observed activity of these compounds is in no way predictated on this mechanistic interpretation. A discussion of the background of this art is given in the commonly assigned application, Ser. No. 330,042, filed Feb. 6, 1973 now U.S. Pat. No. 3,872,171 the disclosure of which is incorporated here by reference."}
{"text": "1. Field of the Invention\nThe present invention relates to a gate driver for driving a thyristor, such as a GCT (gate commutated turn-off) thyristor, a GTO (gate turn-off) thyristor, a static inductive thyristor (SITH), or a power transistor, in a stable condition.\n2. Prior Art\nGCT and GTO have similar turn-on and on-state gate current requirements. The gate driver provides a turn-on pulse with high dI/dt and peak current, and thereafter provides a steady state DC on-state sustaining current. The gate driver sometimes receives negative gate voltage during on-state of the GCT/GTO.\nA new situation, unlike in the case of the GTO, is introduced for driving the GCT. The gate driver for the GCT has its feature in the amount of turn-on peak current and its shape. The turn-on current increase (dI/dt) for a GCT is 10 times to 50 times higher than that of typical GTOs, and trigger peak currents also are selected 5 to 10 times higher in order to reduce turn-on loss. A 6 inch GCT thus may require dI/dt=500A/us and a peak current of 300A.\nWhile such a sharp high pulse safely turns on the GCT accompanying a high rate of rise of anode current, a long pulse duration must be provided after the sharp high pulse to maintain the GCT in on-state, if GCT operation accompanies a small rate of rise of anode current. Thereafter, on-state gate current similar to the GTO\"\"s (approximately 10A for a 4 inch or 6 inch device) is to be applied.\nIf no protection circuit, such as a current limiter, is provided, the gate driver may be damaged when the GCT\"\"s or GTO\"\"s load current changes its flow direction from a forward direction to a reverse direction. When the load current flows in the reverse direction, the load current will flow through a freewheel diode provided in parallel to the GCT. Thus, the GCT\"\"s anode turns negative with respect to its cathode. A parasitic diode in the GCT allows a negative potential to appear at the gate of the GCT.\nThereafter when the load current changes its flow direction again from the reverse direction to the forward direction, the load current may flow again through the GCT. Then, a safe GCT on-condition has to be established again.\nBasically such requirements also are well known from GTO operation. GCT circuits are designed for higher switching frequency, and, as a consequence, freewheel diodes must respond to a considerably higher forward voltage drop, resulting in GCT gate voltage as low as xe2x88x925V during normal converter operation.\nThe peak current is created from a voltage source by a resistor and capacitor circuit. In order to produce a high shooting pulse with a steep and rapid rising (dI/dt=500A/us), and a relatively long down slope pulse (10 us to 40 us) after the shooting pulse, requires several RC combinations with complicated adjustments. In such a design the total loss may exceed 50W for a 6 inch GCT drive.\nA German Patent Laid-open publication No. DE3709150 and a PCT International Publication No. WO9407309 disclose the GTO driver using switched current sources. Inductors are fed to create the sources, and excessive energy is fed back to the power supply.\nBasically in such way, very low loss can be achieved. But four high current switching devices and three diodes are required for generating a turn-on pulse and an on-state current. And the high turn-off peak current has to pass through a series connection of a switching element and an additional diode. Such a circuit cannot be used for a gate driver of GCT.\nU.S. Pat. No. 4,791,350 discloses a gate driver which uses a switch-mode step down current regulator as a source for the GTO\"\"s steady state gate current. The regulator incorporates a switch and a freewheel diode, it\"\"s output is directly coupled to the GTO\"\"s gate, and a high current pulse is generated by a separate circuit having a switch and a resistor.\nIn this way, U.S. Pat. No. 4,791,350 suggests reducing gate driver losses. But if a negative GCT gate voltage should appear, the regulator\"\"s output current will increase without control. Moreover considerable losses are generated in the regulator\"\"s freewheel diode, when gate currents exceeding 5A are required. Also, high losses will result from the high current pulse resistor.\nJapanese Patent laid-open publication No. H3-97315 and EP Patent publication No. 0 416 933 discloses a circuit to solve the problem with negative gate potential. The freewheel diode is connected to the negative supply line. The inductor is charged by the positive supply source. Upon freewheel, the charge in the inductor will be discharged to the negative supply. In this way the circuit can operate stably under all positive and negative GTO gate voltage conditions.\nSuch circuit is applicable for the generation of a small gate current. With high gate current, however, a lot of energy is transferred from the positive supply source to the negative supply source, and it must be transferred back to the positive supply by an appropriate power returning circuit. In the case of GCT, for example, a gate current Ig may be about 10A and a gate voltage Vg is about xe2x88x9220V. Then, a circulating power would be approximately 200W. In contrast to this, the active gate power according to this reference may be as small as Vgxc3x97Ig 0.6Vxc3x9710A=6W. As a result, a tremendous over design of gate current generator and power supply is required, and a loss (approximately 20W-40W) is generated even with high-performance switch-mode circuits.\nEP Patent publication 893883 discloses a gate driver which handles the GCT\"\"s freewheel situation in another way. A bipolar transistor, implemented as emitter follower, is to limit the GCT\"\"s gate current at negative gate voltage. The current is generated from voltage pulses with high efficiency.\nFor gate current Ig up to 2A, a circuit can be designed with appropriate components. At a higher gate current, the bipolar transistor gain will decrease below 20. Then, high base current is required, and loss will increase due to higher VCEsat (saturation of voltage Vce). At Ig=10A, Ib greater than 0.5A will be required, and VCEsat will amount to approximately 1V with the known PNP transistors.\nMoreover, a high base current must be maintained with a GCT gate voltage being maintained not lower than 0.4V. Then, for higher GCT gate voltage and for open circuit condition (no GCT installed), a severe trade-off must be handled, complicating design and limiting the performance of the gate driver.\nA circuit must be found which is appropriate to realize high current GCT drive. It has to generate a sharp, high trigger current pulse for GCT turn-on, a long trigger current tail and high gate current with low loss, and it has to be safe under all gate voltage operating conditions.\nMoreover the circuit has to be such that it can be realized mainly with SMD components and technology in order to allow very compact and cost effective solutions.\nOne aspect of the invention includes a gate driver for driving a thyristor having an anode, a cathode, and a gate by providing a gate current to the gate of the thyristor during an on-command signal, the gate driver comprising:\na turn-on pulse generator for generating a turn-on pulse in response to a leading edge of the on-command signal;\na down converter for producing a down slope current immediately following the firing pulse; and\na current limiter having a MOSFET connected to the gate of the thyristor for supplying current from the down converter to the gate of the thyristor, the current limiter monitoring the gate voltage at the gate of the thyristor and increasing an internal resistance of the MOSFET relative to the negative voltage increase of the gate voltage.\nAccording to the invention, the down converter comprises:\na pattern generator for generating a pattern of a current to be produced from the down converter;\na conduction pulse generator for generating conduction pulses in accordance with the pattern;\na switching element for conducting a current from a power source in response to the conduction pulses and producing a pulse current; and\nan inductor for smoothing the pulse current;\nthe pattern generator generating a pattern having a rising edge and down slope portion after the rising edge, so that the down slope current produced from the down converter decreases relative to the pattern.\nIn a gate driver, the down converter further comprises:\nanother conduction pulse generator for generating other conduction pulses in accordance with the pattern;\nanother switching element for conducting a current from the power source in response to the other conduction pulses and producing another pulse current; and\nanother inductor for smoothing the other another pulse current so that the down slope current is increased.\nIn a gate driver, the down converter further comprises:\nan additional inductor in parallel with the inductor; and\nan additional switching element in series with the additional inductor to produce a greater down slope current.\nIn a gate driver, the down converter further comprises:\nan additional inductor in parallel with the inductor; and\na saturating reactor element in series with the additional inductor to produce a greater down slope current.\nIn a gate driver, the current limiter comprises:\na comparator for comparing an output voltage at the output of the down converter with a predetermined voltage and producing an adjust signal relative to a difference between the output voltage and the predetermined voltage when the output voltage falls below the predetermined voltage, and\nthe MOSFET receiving the adjust signal to change its internal resistance relative to the adjust signal.\nA gate driver according to the invention has as the comparator an operational amplifier.\nIn a gate driver, the comparator comprises bipolar transistors.\nA gate driver according to the invention includes in the current limiter a freewheel diode connected parallel to the MOSFET.\nIn the gate driver, the turn-on pulse generator comprises a capacitor, a diode connected parallel to the capacitor, a switching element, and a reactor to produce a sharp high pulse.\nThe gate driver may comprise a bias current generator to provide a low level bias current to the gate of the thyristor.\nA gate driver according to the invention is a GCT."}
{"text": "1. Field of the Invention\nThis invention relates to a nozzle or head of the type mounted on the ends of flexible hoses and used for the cleaning of pipes, ducts and waste water conduits using high velocity streams of liquid, and more particularly to such a nozzle which produces a pulsating flow.\n2. Description of the Prior Art\nWaste conduits containing caked or loosely adherent sludge, sediment, stone and other solid articles have been cleaned with the aid of a self-propelled nozzle or head affixed to the end of a flexible hose and having a plurality of rearwardly directed jets of high velocity fluid, usually water. In these systems, the reaction force upon the nozzle propels it through the conduit while entraining the attached length of flexible hose therealong. The fluid jets concurrently loosen adherent materials and carry them back through the conduit. When the nozzle is retracted in the direction of the advance of the nozzle by drawing upon the hose or a cable affixed to the nozzle, the rearwardly directed streams of fluid act to scoop sediment and contaminants away from the walls of the waste conduit and thereby provide an additional loosening and cleansing effect. Examples of prior art self-propelled nozzles are shown in U.S. Pat. Nos. 1,176,518 issued to Burns, 1,628,070 issued to Sladden, 3,080,265 issued to Maasberg, and 3,321,184, issued to Goss.\nThe flow rate of fluid emitted by the nozzle and the pressure of the fluid are limited by the output of the pump supplying fluid to the nozzle and the capacity of the flexible hose which connects the pump to the nozzle. The pressure and flow rate of the fluid are regulated in the nozzle by the size and number of the jets or orifices in the nozzle through which the fluid flows. With the pump producing a maximum flow rate at maximum pressure of the pump outlet, the pressure and flow rate of fluid through the nozzle are limited by the capacities of the pump and the connecting hose.\nWhile the cleaning effect produced by these prior art nozzles has been sufficient where the degree of sedimentation is not excessive and where the walls of the conduit are coated with a loosely adherent layer of sludge, these nozzles have encountered difficulty in removing heavy accumulations of adherent sludge or caked-on material even when the water dispensed by the nozzle is at high velocity with the maximum pressure and volume possible.\nIt has been proposed to produce a pulsating jet of fluid from the nozzle in order to increase the effectiveness of the nozzle in removing caked sludge from the walls of the conduit. Such a pulsating jet has been produced by using a pump which produces a pulsating flow of fluid through the flexible hose to the nozzle. This pulsating system requires a specially adapted pump which can be quite expensive. In addition, the effectiveness of this pulsating system has been limited by the capacity of the pump and the capacity of the flexible hose to carry a given maximum rate of fluid through the hose from the pump to the nozzle. Thus, when the pulsating flow from the pump reaches its maximum flow rate, this maximum flow rate cannot be any greater than the maximum flow rate produced by a continuous flow because it reaches the maximum capacity of the hose. The effectiveness of this pulsating spray is, therefore, not much greater than a continuous spray. Another problem with this pulsating system is that the pump is often a considerable distance away from the nozzle and this pulsating flow produced by the pump loses some of its effectiveness as the fluid travels through the flexible hose to the nozzle."}
{"text": "1. Field of the Invention\nThe present invention relates to an immersion microscope objective for use in a laser scanning microscope that uses multi-photon excitation.\n2. Description of the Related Art\nA known microscopic fluorescence observation means uses multiphoton excitation in fluorescence observation. During multiphoton excitation, a fluorescent object is illuminated with a light beam having wavelengths of integral multiples of an inherent absorption wavelength so as to cause excitation nearly equivalent to that caused by light having a wavelength equal to the inherent absorption wavelength. Multiphoton excitation is a non-linear phenomenon. For example, in the case of two-photon excitation, the excitation occurs at a probability proportional to the square of the intensity of the excitation light.\nOn the other hand, excitation light collected by a microscope objective lens has a light intensity that is inversely proportional to the square of the distance from the focal plane. Microscope multiphoton excitation occurs only in the vicinity of the focal point and thus fluorescence is emitted only from the focal plane.\nBecause of the above, a confocal pinhole, which is provided in the detection system of a common confocal microscope in order to shield the detector from fluorescence emitted anywhere other than at the focal plane, is unnecessary when using a multiphoton excitation laser scanning microscope. Advantageously, fluorescence in samples is subject to less discoloration because the excitation occurs only at the focal plane.\nIn multiphoton excitation, infrared light (that has longer wavelengths than light ordinarily used for excitation) is generally used as the excitation light. Generally speaking, light having longer wavelengths tends to scatter less (i.e., Rayleigh scattering). On the other hand, infrared excitation light reaches deep inside scattering samples such as living samples. This allows the deep parts of a living body to be made observable, as those parts are otherwise impossible to observe using only visible light. Furthermore, infrared light is less photo toxic than ultraviolet and visible light. This allows for observation of living samples with minimum damage.\nFor example, serotonin (a substance found in the brain) is self fluorescent, having absorption wavelengths in the ultraviolet range. However, ultraviolet light is highly photo toxic and does not reach deep inside the brain. In such circumstances, multiphoton excitation laser scanning microscopes work effectively. The ultraviolet light and infrared light differ in wavelength by a factor of approximately three. Thus, the infrared light can be used to excite serotonin by three-photon excitation.\nAs described above, multiphoton excitation fluorescence observation has significant merits and currently is a very effective type of microscopic observation.\nHowever, microscopic observation using multiphoton excitation has some technical difficulties. For example, for multiphoton excitation to occur, two or more photons at a time must collide with a molecule of the fluorescent object and be absorbed. In order for this to occur, a significantly high photon density must be created at the focal point of the microscope objective lens. In other words, a microscope objective lens having a large numerical aperture with aberrations that are properly corrected must be used. As the excitation light is infrared light, it is important that aberrations of the microscope objective lens be favorably corrected for the infrared light.\nFor observation of the deep parts of a sample, aberrations resulting from the refractive index of the sample cannot be ignored. In order to avoid a deterioration of the fluorescence efficiency due to such aberrations, it is desirable that a correction collar be provided that corrects for aberrations that vary according to the depth of observation in the sample.\nOn the other hand, even though the excitation light is strong, only a small intensity of fluorescent light emerges from a sample. Therefore, the emerging fluorescent light has to be detected with a minimum loss. Thus, the microscope objective lens should be composed of a minimum number of lens elements and lens groups so as to optimize the optical path to the detector.\nFurthermore, ultraviolet light should be detected in the case of three-photon excitation. However, conventional optical glass has a low transmittance to ultraviolet light. Therefore, the microscope objective lens should be constituted by specific types of optical glass that have a high transmittance to ultraviolet light.\nA microscope that uses multiphoton excitation is often used in conjunction with a technique known as patch clamping that is widely used in cellular biology studies. When used in this manner, a sufficient working distance has to be reserved between the leading end of the immersion microscope objective and the sample. In other words, a large working distance and a sufficient access angle at the leading end of the immersion microscope objective need to be provided.\nThe excitation light is infrared light and the emerging fluorescence lies within the visible and ultraviolet ranges. Therefore, the fluorescence to be detected is likely to be subject to Rayleigh scattering by the sample. Consequently, light emerging from one point may become diffused before it enters the immersion microscope objective. Therefore, it is desirable that the immersion microscope objective be provided with a large field of view so as to collect the scattered fluorescence without loss.\nWhen the field of view of the immersion microscope objective is extended, the pupil diameter of the immersion microscope objective increases proportional to it. Even if it is assumed that an immersion microscope objective having a very large input pupil diameter is developed, the performance cannot be displayed unless the input laser beams of the confocal microscope device fill the pupil diameter in the input pupil position of the immersion microscope objective. Therefore, Japanese Laid-Open Patent Application No. 2008-040154 proposes a technique for adjusting the input beam diameter in such a way as to approximately match the pupil diameter of the immersion microscope objective with an input beam diameter. According to this technique, the confocal microscope device can adjust the input beam diameter in such a way as to approximately coincide with the pupil diameter of the immersion microscope objective. However, since the input beam diameter is adjusted before beams enter a laser beam deflection means (a two-dimensional scanning means), the maximum input beam diameter is defined by the laser deflection means.\nIn order to expand this, it is necessary to increase the device size of the laser beam deflection means such as a galvano-mirror and the like. However, the larger the size of the galvano-mirror becomes, the larger a deflection angle for scanning the specimen range becomes to reduce the scanning speed. In a resonant galvano-mirror provided with two galvano-mirrors near a position conjugate with the aperture position and capable of two-dimensional scanning, the distance must be extended in order to avoid the interference between the two galvano-mirrors. Thus, it deviates from the ideal position to increase illumination unevenness.\nTherefore, it is important to combine an optimum immersion microscope objective with the maximum input beam diameter defined by the laser beam deflection means of the confocal microscope device."}
{"text": "Mobile devices such as wireless devices, including, for example, cellular telephones, smart phones, laptop computers, notebook computers, tablet devices (e.g., iPad by Apple®) are ubiquitous in modern society. Use of such mobile devices while operating a vehicle, however, can be hazardous. The problem is exacerbated for inexperienced operators of the vehicle, such as youngsters just learning how to drive. Rates of vehicular accidents where mobile devices are involved are rising, especially with teenagers. Text messaging while operating a moving vehicle can be dangerous and has been linked with causing accidents. More generally, operating any keyboard or other interactive device while operating a vehicle can be dangerous.\nThus, the widespread adoption of mobile devices and common use of the devices while driving has raised concerns about the distraction of drivers. A driver speaking, text messaging, or using a software application on a mobile telephone may become mentally distracted from driving and lose control of the vehicle that he or she is driving. Thus, it is not uncommon to see an individual involved in an accident who was speaking or text messaging on a mobile device rather than paying attention to the road. Studies now suggest that individuals speaking on mobile telephones while driving a car may be as impaired as a person who drives while intoxicated. Not only is the driver mentally distracted, but eyes of the driver are diverted for dialing, looking to see who an incoming call is from.\nIt would be highly desirable to detect the presence of a mobile device, such as a wireless device, within a vehicle and control or inhibit the operation of the mobile device."}
{"text": "Atherosclerosis, the clogging of arteries with plaque, is often a cause of coronary heart disease or vascular problems in other regions of the body. Plaque is made up of fat, cholesterol, calcium, and other substances found in the blood. Over time, the plaque hardens and narrows the arteries. This limits the flow of oxygen-rich blood to organs and other parts of the body.\nBlood flow through the peripheral arteries (e.g., carotid, iliac, femoral, renal etc.), can be similarly affected by the development of atherosclerotic blockages. Peripheral artery disease (PAD) can be serious because without adequate blood flow, the kidneys, legs, arms, and feet may suffer irreversible damage. Left untreated, the tissue can die or harbor infection.\nOne method of removing or reducing such blockages in blood vessels is known as rotational atherectomy. In some implementations, a drive shaft carrying an abrasive burr or other abrasive surface (e.g., formed from diamond grit or diamond particles) rotates at a high speed within the vessel, and the clinician operator slowly advances the atherectomy device distally so that the abrasive burr scrapes against the occluding lesion and disintegrates it, reducing the occlusion and improving the blood flow through the vessel."}
{"text": "1. Field of the Invention\nThe present invention belongs to a technical field related to a display device in which a semiconductor device (typically, a transistor) is used as a device, in particular, a light emitting device represented by an electroluminescence display device, and to a technical field related to an electronic apparatus equipped with the display device in an image display portion.\n2. Description of the Related Art\nIn recent years, development of liquid crystal display devices and electroluminescence display devices in which transistors (particularly, thin film transistors) are integrated on a substrate have progressed. Such display devices are respectively characterized in that the transistors are formed on a glass substrate by using a thin film formation technique, and the transistors thus formed are disposed in pixels each arranged in matrix and made to function as a display device for image display.\nA variety of specifications are conceivable, which are required for areas (hereinafter, referred to as pixel portions), in which the image display is performed in the display device. However, the following are given as examples thereof: a large number of dots and high definition are ensured; an area of an effective display region in each pixel is large and bright image display is possible; and the pixel portion involves no defects that may induce point defects or line defects in its entirety. In order to achieve those specifications, not only the performance of the transistors arranged in each pixel should be satisfactory but also a technique of forming the transistors while increasing yield through a stable process is necessary.\nFurther, in an organic electroluminescence display device among the electroluminescence display devices, an organic compound is used for a light emitting element serving as a light emitting source. Accordingly, a measure for suppressing deterioration of the organic compound is most highly required in ensuring its reliability. In other words, in order to achieve a highly reliable display device, attention must be paid not only to an influence of an accumulated damage in the process during manufacturing the device but also to the subsequent deterioration with time, which results from the accumulated damage.\nIn the above-mentioned circumstances of development, the applicants of the present invention are most concerned, in the present conditions, with problems such as variation and shift of a threshold voltage in the transistors, which arise due to the accumulation of plasma damages on an insulating film etc. in an etching process."}
{"text": "1. Field of the Invention\nThis invention relates generally to grips or handles for office, industrial or household devices, which have an anti-microbial treatment so that the transfer of germs and bacterium from one handler to another is reduced or eliminated.\n2. Description of the Related Art\nIt is very common for people to share hand held products such as office products, industrial and household tools. This is especially true in manufacturing facilities where multiple workers share the same tools. For instance, in a circuit board fabrication facility, soldering irons for soldering or desoldering components to circuit boards arc often shared by workers in subsequent shifts. One of the problems of sharing such tools is that germs and bacterium may pass from the hands of one person to the tool and from the tool to another person. This may lead to workers becoming sick and unable to work, which lowers the productivity of the workers. Therefore, there is a need for reducing the chance of spreading germs through people's hand as they share things such as office or industrial products and tools."}
{"text": "Cartridge case trimmers, as contemplated herein, are of the mini-lathe variety, i.e. a selected cartridge case is locked in place and a blade is mounted on, or is an integral part of, a rotating shaft that is designed for controlled movement toward the mouth of the cartridge to engage and trim the case back to a desired length. Trimmers such as described above are used for refurbishing cartridge cases, i.e. where the cartridge has been fired and the spent cartridge case is cleaned, reshaped and reloaded. Such reloading of cartridge cases is popular among users of large numbers of cartridges, e.g., hunters, target shooters etc. There is a substantial savings to the reloader and he can, in effect, customize his own cartridges.\nThe trimming operation is made necessary because the explosion that takes place in the cartridge causes some expansion and elongation of the material of the cartridge case. The cartridges are designed to precisely fit the chamber of a specific type of firearm and thus the need for the trimmer.\nThe problem that is addressed by the present invention concerns the necessity of trimming cartridge cases to different lengths, i.e., different caliber cartridges (for different firearms) have different lengths. The enthusiast who loads his own cartridges is very likely to have a number of firearms and thus a number of different caliber cartridge cases that must be accommodated by the trimmer.\nThe trimmers presently in common use are equipped with an adjustable stop member that can be set so that the movement of the cutter into the cartridge case mouth is stopped at precisely the desired length of the cartridge. However, every time a different cartridge is to be trimmed, the reloader must readjust the stop member. Readjustment is a time consuming and tedious trial and error process."}
{"text": "1. Field of the Invention\nThe present invention relates in general to a bumper structure of a motor vehicle, and more particularly to a bumper structure including a bumper fascia.\n2. Description of the Prior Art\nSome of light duty motor vehicles are equipped with bumper structures of a type having a plastic bumper fascia.\nAn example of prior art bumper structure of the above described type is shown in FIGS. 4-7 of the accompanying drawings. As is shown in FIG. 4, some of conventional bumper fascias 10 are integrally molded so as to have openings 12 and 14 into which head lamps and a radiater grille are fitted. As is seen from FIG. 5, the bumper fascia 10 is fastened, at its rear end vertical flanges 16, to front end vertical flanges 22 of a front fender panel 24 through bolts 18, nuts 20 and retainers 26.\nHowever, the bumper structures of he above-mentioned type has the following drawbacks.\nAs the size of the bumper fascia 10 becomes increased, the same tends to be distorted in shape due to inevitable thermal contraction at its molding process.\nThus, if the bumper fascia 10 having been outwardly distorted is mounted onto the front fender panel 24, it tends to occur that, as is seen from FIG. 7, a rear upper end of the bumper fascia 10 mismates with a front upper end of the front fender panel 24 leaving a considerable gap \"A\" therebetween. Designated by numeral 28 in FIG. 7 is a hood panel for the engine room. Furthermore, as is seen from FIG. 6, because of the thermal contraction of the bumper fascia 10, it tends to have an undesirable gap \"B\" betwen a rear outer end of the bumper fascia 10 and a front outer end of the front fender panel 24. These gaps \"A\" and \"B\" lower the external appearance of the front face of the motor vehicle."}
{"text": "Machines such as wheeled compactors, loaders, trucks, and other mobile equipment are used to perform many tasks. To effectively perform these tasks, the machines require an engine that provides significant torque through a drive system to one or more ground engaging traction devices (e.g., wheels, tracks, etc.). Drive systems for such machines usually include either a mechanical power-shift transmission that is connected to the engine by way of a mechanical or hydraulic torque converter, or a hydrostatic transmission (hystat) that pressurizes fluid with an engine-driven pump for powering a motor connected to the traction device. Recently, more attention has been given to the use of hystats because hystats are known to achieve higher fuel efficiency and extended functionality when compared with a power-shift transmission. A hystat provides an infinitely variable torque-to-speed output ratio within its overall range through the pairing of a variable displacement pump and a fixed- or variable-displacement motor.\nAlthough higher fuel efficiency and extended functionality can be achieved with hystats, many operators still do not prefer machines equipped with hystats because hystat systems respond to operator commands (e.g., throttle commands, directional shits, shifts to neutral, work tool commands etc.) with different characteristics (e.g., acceleration rates, deceleration rates, rimpull torque, etc.) than the familiar and comfortable characteristics of power-shift systems. For instance, when an operator releases the acceleration pedal or throttle of a machine equipped with a hystat, displacements of the pump and/or motor are immediately neutralized and the machine either continues traveling at about the same speed for an extended period of time (when both the pump and motor are neutralized) or quickly stops (when only one of the pump and motor are neutralized). Therefore, operators must modulate the acceleration lever or pedal to move through intermediate displacement positions in a gradual manner in order to travel at an increasingly slower speed. This can prove very difficult in rough terrain common to many worksites, especially for unskilled operators, and can cause undue operator fatigue.\nOne attempt to vary the response characteristics of a hystat drive system is described in U.S. Patent Publication No. 2013/0104532 (the '532 publication) by Ries et al. that published on May 2, 2013. The '532 publication describes a hystat drive system having a transmission input device that allows an operator to select a transmission setting for operating the machine with a desired virtual gear or range of speed-to-torque ratios. Based on the operator's selection, a controller selects a relationship from an acceleration map that corresponds with the transmission setting. The controller then uses the relationship to adjust the displacement of a pump and/or a motor in the hystat system during an acceleration process. When the operator initiates a coasting or braking operation, the controller selects a coasting deceleration map or a braking deceleration map, respectively, and adjusts the pump and/or motor displacements to slow the vehicle by a desired amount over a desired period of time after initiation of the coasting or braking operation. The desired period of time is a function of the machine's travel speed when the coasting or braking operation was initiated.\nAlthough the system of the '532 publication may be somewhat effective at changing the acceleration rate and deceleration rate of a machine equipped with a hystat, it may not be optimal. In particular, because the speed change characteristics of the '532 publication are based on a gear selection, acceleration and deceleration characteristics of the machine may not be optimal under varying conditions. Further, the single control mode of the '532 publication may not allow the machine to be effectively controlled by operators of varying skill and/or familiarity levels.\nThe system of the present disclosure solves one or more of the problems set forth above and/or other problems of the prior art."}
{"text": "Lamps for vehicles are known, designed so that the light emitted or reflected by them to make it possible to recognize a symbol or text.\nPatent DE 202005003397 describes a rear lamp for motor vehicles or bicycles provided with a customized reflector designed to highlight its reflections as an illuminated symbol, and in the absence of a light source, the symbol remains visible as a dark shape.\nPatent DE 202006008944 relates to a motor vehicle lamp which includes a logo or symbol, evidenced by the intense illumination when the driver of the vehicle brakes or signals change of direction, an effect obtained by using extra-bright LED inside the lamp.\nThe disadvantages of known solutions is that they require the manufacturing of the complete lamp assembly according to customization requested by the beneficiary, requiring significant resources to prepare the launch of the production; another disadvantage is that personalization is highlighted only at night, in the presence of incident light sources or certain conditions: braking, signaling change direction and so on, without personalizing be easily observed at all times."}
{"text": "Additive manufacturing (AM) techniques such as fused deposition Modeling® (FDM®) and fused filament fabrication (FFF) are rapidly gaining popularity in a variety of manufacturing industries due to AM's ability to create unique geometries and process unique material compositions. In particular, AM can operate on a variety of length scales with a large number of thermoplastic feedstock materials.\nAn additive manufacturing process includes building parts from selectively dispensing material through a nozzle or orifice. Multiple layers are built on top of each other to create a part or solid object. Material such as polymer is delivered to an extruder as a filament coiled in a spool, and is then melted during extrusion. That is, the polymer, typically held in the machine as spooled filament, is fed by the extruder into the liquefier where it is heated and melted, by receiving heat by conduction with heated walls in the liquefier. The extruder therefore feeds the filament through the liquefier, such that molten polymer is dispensed through the nozzle at the end of the liquefier chamber, which reduces the diameter of the molten polymer at its output. The extruder and liquefier assembly are directed in the X-Y plane using a gantry to build a cross-section of the part. Parts are built on a movable build platform within a chamber designed to control the printing environment. When printing in thermoplastic, the platform and/or its chamber are typically heated to reduce residual stresses, within the part, that arise during cooling after the deposition from the nozzle.\nTwo important components during the additive manufacturing process are the filament extrusion mechanism and the filament liquefier. Typically, a pinch-wheel mechanism that includes a drive wheel and a pinch or knurled wheel form a filament extrusion mechanism that drives the filament linearly by applying force along the side of the filament. Traction between the drive wheel and filament is provided by the clamping force applied by the pinch wheel. Driving the filament in this way generally causes the filament to be compressed between itself and the drive wheel, particularly when the drive wheel employs knurls, which dig into the filament. A filament fails by shear when an amount of force applied to the filament by the drive wheel exceeds the average shear stress over the area of engagement with the filament. Generally, the filament liquefier heats the polymer by conducting heat energy to the polymer from the liquefier wall and melts the polymer within the cavity before reaching the extrusion nozzle where its cross-section is reduced.\nThe production of parts using additive manufacturing is limited, for instance, by the design of current extruders, liquefiers, and gantries. For example, traditional extruders inhibit extrusion rates by providing a filament engagement area that limits the force that can be applied or transmitted by the extruder on the filament; traditional liquefiers do not provide an adequate heat rate necessary for expedient extrusion; and traditional gantries are speed and acceleration limited and cannot move with the agility needed to match high-throughput extrusion rates.\nThat is, in general, extrusion force in current extruders is limited by the filament shear area. Extrusion failure occurs when the extruder stress is greater than the shear stress of the filament shear area but less than the pressure drop caused by the liquefier.\nMoreover, current liquefiers are limited by the maximum heat rate into the filament during extrusion, because they generally rely on conduction heat transfer from the heated liquefier walls. Due to the poor thermal conductivity of the filament and the conduction heat transfer rate, filament feed rate must be kept slow to allow full heating of the filament to the required melt temperature. Thus, with increasing extrusion rates, the filament receives less heat from the conduction liquefier. If the filament is not heated to the required temperature, the filament does not melt in the liquefier and instead jams in the nozzle, thereby preventing extrusion.\nCurrent gantries for positioning a nozzle and build platform are effective at low dynamic performance, but traditional stepping motors used with gantries are speed limited due to fundamental operating principles of the actuators and their open-loop control scheme. On the other hand, simply using larger actuators causes an increase in mass and inertia of the gantry, thereby slowing the gantry response time despite having a higher torque capacity.\nAs a result of the above limitations of current additive manufacturing techniques, including insufficiently heating the filament, a high pressure needs to be exerted on the liquefier in order to properly extrude the filament. However, the high pressure or force on the filament may cause filament shear failure.\nAccordingly, there is a need for systems, devices, and methods for high-throughput additive manufacturing that, among other things, are designed to increase the maximum force capable by the extruder, the heat rate into the volume of the filament, and the rate of deposition of the material of the filament. Moreover, there is a need for such systems, devices, and methods to be applicable to polymer resins and high-performance thermoplastic, and adaptable to large scale extrusion AM systems, in a variety of sectors such as medicine, product design, cinematography, education, aerospace, and advanced materials, among other sectors."}
{"text": "The popularity of the Internet, coupled with the increasing capabilities of personal/mobile electronic devices, has provided consumers with the ability to enjoy multimedia content almost anytime and anywhere. For example, live (e.g., sports events) and video on demand (VOD) content (e.g., television shows and movies) can be streamed via the Internet to personal electronic devices (e.g., computers, mobile phones, and Internet-enabled televisions).\nIn a client-server system, a server may provide a multimedia stream to each client device that requests the stream. However, providing multiple copies of the stream may exhaust available bandwidth at the server. To alleviate the demand on the server, the stream may be syndicated to multiple servers and made available from any of the servers. A load balancing system may be utilized to determine which particular server should provide a stream to a particular destination device."}
{"text": "Bridges, buildings, railroads, pipes, vessels, tanks, and other metal or steel welded structures are a vital part of modern infrastructure. These structures are typically fabricated from sets of discrete metal sub-components that are welded together to form a critical component and/or a total system. Welded seams and other weld joining points must be fused into a welded unit having satisfactory strength to ensure building code compliance and to achieve proper structural integrity of the entire system, and/or to ensure meeting the purpose of the design for which it was intended. For example, a pressure vessel or fluid tank must have water tight exterior as well as provide structural support for the entire vessel or tank system. Pipes and vessels similarly must have water or gas sealed, welded seams to ensure the integrity of the pipe/vessel and to properly isolate the fluid or gas held by the pipe from the environment. In addition, these systems deteriorate over time due to operational and environmental factors such as, residual and applied stresses, vibration, rain, snow, strong winds, temperature variance, earthquakes, oxidation, material fatigue, and other changes that occur over the passage of time. Hence, nondestructive testing of welds and metal structures, and their components, are utilized after initial fabrication, installation, and periodically thereafter, to ensure a structure's integrity. Further, both for new construction and routine periodic maintenance, careful analysis of the weld joints in each structure is necessary to ensure satisfaction of various weld specifications, industry codes, and construction regulations. For example, the American Society Mechanical Engineers (“ASME”), the American Welding Society (“AWS”), and the American Petroleum Institute (“API”), among others, each have their own welding codes, procedures, and specifications.\nModern inspectors use non-destructive test (“NDT”) equipment to inspect constructed metal (e.g. steel) structures and their weld joints. These inspection devices use ultrasonic wave generators to take digital “snap-shots” of welds from which an inspector may verify weld integrity and to ensure compliance with welding codes and specifications. Ultrasonic technology is used to detect internal and surface breaking flaws in the weld and the base metal, which are not visible externally, and is based on the principle that a gap or defect in the weld changes the propagation of ultrasonic sound through the metal. One common method of NDT testing uses conventional, single-probe ultrasonic testing requiring an operator's interpretation of a screen similar to an oscilloscope screen that presents time and amplitude information. Another method uses an array of ultrasonic phased array sensors to test a structure. Such methods can be combined into a single digital piece of inspection equipment that uses phased array (“PA”) and time-of-flight (“TOFD”) diffraction methodologies to provide a three dimensional image of a weld displayed on a color screen. An inspector then evaluates the potential for a flaw or defect in the weld by reviewing the screen. Such ultrasonic testing (“UT”) equipment is typically highly mobile, and allows for the recordation of ultrasonic data for the analysis of welded areas in joined metal pieces. For example, Olympus NDT, Inc. markets and sells ultrasonic units through its OmniScan™ and Epoch™ lines of weld flaw detectors. The Olympus ultrasonic inspection systems include conventional ultrasonic flaw detectors, which use ultrasonic waveforms to detect flaws, and advanced ultrasonic phased array flaw detectors, which create internal cross-section images of the areas being inspected. These ultrasonic flaw detectors may also be configured to do a phased array ultrasonic testing (“PAUT”) inspection that produces encoded digital data points on welds which may be further processed at a later time. The data may also be reviewed later at a time convenient for an inspector and in an environment away from the component or structure's location, which is typically more suitable for detailed analysis work. Such ultrasonic data is recorded and saved in large data files which may be retrieved for evaluation by an inspector using specialized software applications. While these data files are readable by the ultrasonic testing device, software applications also exist that assist in the evaluation and visual display of such inspection data on common computing devices, such as a PC. An example of such applications is the OmniPC™ analysis software also available from Olympus NDT, Inc., along with an additional analysis tool set called TomoView™. Both of these tools allow for a more precise and reliable review of the three dimensional PA data by an inspector.\nA phased array data file consists of captured data representative of continuous A-scans along a weld which may be processed to create a three dimensional data set representative of the top, side, and end views of a weld joint. Within each A-scan, a series of recorded data points record an intensity or amplitude value from 0% to 100% of ultrasonic signal reflections. These data points are then saved as OPD, RTD or TV file formats for subsequent analysis. Since the recordation of ultrasonic data is correlated to the exact PA probe position on the structure being tested, the inspector may use the analysis software including various sets of sophisticated analysis tools to review the recorded data in a number of geometric views and orientations (e.g. top view, side view, end view) to improve inspection review accuracy.\nFurther information regarding the use of phased array UT equipment, configuring such UT equipment, establishing a test scan plan applicable to a particular inspection or weld joint design situation, the recording of that data in various file formats, the physics and geometries of the ultrasonic sound beams and resulting scanning views in UT, the use and applicability of “data” libraries, the storing of testing data files, the usage of different types of scan views, the visual analysis of weld flaw indications, and the generation of inspection reports based upon UT shall not be discussed herein as such information is well known known in the NDT industry and not necessary for a complete understanding of the disclosed invention. However, Applicant references and hereby incorporates by reference the treatise UT Classroom Training Book, Paul T. Marks, ISDN No. 978-1-57117-345-4 (e-book), published by The American Society for Nondestructive Testing, and two treatises published by Olympus NTD, Inc: (1) Introduction to Phased Array Ultrasonic Technology Applications, third printing 2007, ISBN No. 0-9735933-4-2; and (2) Advances in Phased Array Ultrasonic Technology Applications, 2007, ISBN No. 0-9735933-4-2. The books may be obtained at the ASNT website www.asnt.org or the Olympus NDT resources website www.olympus-ims.com. These treatises explain the above subjects in detail and the general theory of UT using modern equipment. Further, Applicant references and hereby incorporates by reference U.S. Pat. Nos. 8,156,813B2, 8,577,629B2, 9,032,802B2, 9,081,490B2. These patents discuss and disclose background information regarding the electronics and theory behind PA ultrasonic testing.\nNevertheless, even with modern PA and time-of-flight UT devices, and even when inspection analysis is conducted in an environment conducive for careful study, the data analysis and reporting process can be a tedious and fatiguing task for inspectors. For example, metal pipe and plate structures typically have girth welds and long seams that must be inspected. PA ultrasonic and time-of-flight, diffraction inspection for those welds seams can produce extremely large data files requiring many hours of data review and analysis of all data points along the weld seams by an inspector. Usually most of the data points are nominal, satisfactory welds, creating a monotonous review period and potentially reducing weld flaw recognition by an inspector due to fatigue. Hence, what would enhance the inspection process would be a system for focusing an inspector's attention on actual weld defect indications, by excluding data points that present satisfactory and acceptable weld characteristics. Such a system would improve an inspector's efficiency and accuracy in conducting new or reviewing prior weld inspections, thereby saving time and money."}
{"text": "Governments of nations from all around the world are currently establishing a host of security programs to ensure the protection of their national treasures, infrastructure, citizenry, and to establish social order. One example of such programs is in the detection of illicit and hazardous materials such as drugs, explosives, chemical and biological contaminants.\nFor the detection of these abovementioned materials, two approaches are normally used. The first approach involves the use of electronics-based detection machines and instruments. These machines and instruments typically operate using methods based on a wide array of spectroscopic, electrical or magnetic principles. From these principles have emerged traditional detection machines based on technologies such as the x-ray, gas chromatographic and nuclear quadruple resonance machines, to emergent detection machines based on chemical and electronic nose nanotechnology based sensors, and non-ionizing terahertz, infrared and fluorescence spectroscopy, interferometry and imaging techniques.\nThe major advantages of detection machines and instruments include their parts-per-billion sensitivity to vapor detection and the fact that different technologies and principles of operation can be used simultaneously within a singular detection configuration.\nHowever, a disadvantages imposed by the use of machines or instrument detection is that the machine itself is not capable of an autonomous search-and-detect process of detection within a defined space or area. It cannot “see.”\nThe second approach to detection is of a biological nature and takes advantage of the olfactory capacity of creatures such as canines, bees, rats and fishes such as sharks and the catfish to detect illicit and hazardous materials, specifically employing their capacity to detect sub-parts-per-million amounts of a scent or odor exuded by a material and then trace the odor back to its source, and the fact that they are easy to train by humans in search-and-detect operations.\nHowever, a major limitation of detector creatures is the fact that they cannot communicate to humans the exact nature of the illicit or hazardous material what they are positively alerting to nor can they confirm the nature of the scent trail they are following. Also, their senses of detection are wholly based on principles pertaining to olfaction.\nIt would be beneficial to have a bio-machine hybrid detector that combines the autonomous search-and-detect capability of a detector creature that is trained in the art of detection of illicit and hazardous materials with the multi-perspective material identification and characterization capability of electronics-based detector machines and instruments, as such a detector will have superb advantages over currently available biological or electronic detection systems in the detection of materials such as explosives, narcotics, laundered currency, fugitives and agricultural produce."}
{"text": "1. Field of the Invention\nThe present invention relates to a system configured to provide a three-dimensional image to a viewer having liquid crystal shutter glasses or the like mounted thereon, a three-dimensional image correction system applicable to a three-dimensional image display system having a function of correcting a pop-up amount, a pull-in amount, and the like of a three-dimensional image to be perceived by a viewer, a three-dimensional image correction method, a three-dimensional image display device, a three-dimensional image reproduction device, a three-dimensional image provision system, a program, and a recording medium. More specifically, the invention relates to a system provided with a correction computation unit configured to correct a dynamic range on the basis of a disparity amount detected from image information for a three-dimensional image composed of an image for a left eye and an image for a right eye and information on a screen size and the like of a display device configured to display a three-dimensional image, in which the disparity amount is corrected in the corrected dynamic range while corresponding to display device information so as to be able to adjust a pop-up amount, a pull-in amount, and the like of a target object while corresponding to the display device information, and also even in a case where a specification of the display device is varied, not only it is possible to realize the pop-up amount, the pull-in amount, and the like preferred by the viewer, but also it is possible to accurately express the pop-up amount, the pull-in amount, and the like intended by an image creator.\n2. Description of the Related Art\nIn recent years, a slimmer television set has been realized with a higher definition and a larger screen. For the next paradigm shift (scientific revolution), researches are actively carried out in the field of three-dimensional display. In a three-dimensional display technology, a method of utilizing a binocular disparity to allow the viewer to perceive a three-dimensional object is generally used, and a large number of achievements have been reported.\nFor example, a three-dimensional image system is developed which is configured to provide a three-dimensional image to the viewer having the liquid crystal shutter glasses or the like mounted thereon. For a three-dimensional image display device in which the viewer has the liquid crystal shutter glasses mounted thereon and views the three-dimensional image, in order to present the three-dimensional image by utilizing the binocular disparity, a display method is adopted in which different images are arranged to enter into the left eye and the right eye.\nAccording to such a display method, a disparity amount is set between the image for the left eye and the image for the right eye while targeting the same object. The disparity amount herein is a difference between a convergence angle on a display surface and a convergence angle on a three-dimensional surface to be perceived based on image information for a three-dimensional image and is, between a pixel displacement amount of a binocular disparity image on the display surface and a standard visual distance from the display surface to eyes of a viewer, substantially approximated by (the pixel displacement amount)/(the standard visual distance). By adjusting this disparity amount, it is possible to set a pop-up amount, a pull-in amount, and the like of the target image. For example, in a case where three-dimensional images are viewed in a cinema or the like, a screen size is fixed. Thus, the pop-up amount, the pull-in amount, and the like of the object which are intended by an image creator and also with which the viewers do not feel unpleasant sensation or sense of discomfort can be substantially uniquely decided.\nIn association with this type of the three-dimensional display technology, Japanese Patent No. 3749227 discloses a three-dimensional image processing method and device. This three-dimensional image processing device includes an instruction obtaining unit, a disparity identifying unit, and a disparity control unit. When a three-dimensional image is displayed on the basis of a plurality of viewpoint images corresponding to different disparities, the instruction obtaining unit obtains a response from a user whether or not the three-dimensional image displayed in various disparities is acceptable. On the basis of the obtained response, the disparity identifying unit identifies an appropriate disparity as the disparity acceptable by the user. On the premise of this, when another three-dimensional image different from the three-dimensional image is displayed, in order that the user can accept the other three-dimensional image, the disparity control unit applies a processing on the other three-dimensional image on the basis of the identified appropriate disparity. With the three-dimensional image processing device having the above-mentioned configuration, it is possible to generate or display the three-dimensional image and the like suitable to human physiology. Also, the three-dimensional sensitivity can be adjusted through a simple operation."}
{"text": "As one of a method of forming a high quality semiconductor film, there is an epitaxial growth technique of forming a single crystal film on a substrate such as a wafer by vapor phase growth. In a vapor phase growth apparatus using the epitaxial growth technique, a wafer is mounted on a support unit in a reaction chamber which is maintained in a normal pressure or a reduced pressure. Next, while heating the wafer, a process gas such as a source gas which is a raw material for film formation is supplied from, for example, a shower plate in an upper portion of the reaction chamber to a surface of the wafer. Thermal reaction of the source gas occurs on the surface of the wafer, so that a single crystal epitaxial film is formed on the surface of the wafer.\nIn recent years, as a material for a light emitting device or a power device, a GaN (gallium nitride)-based semiconductor device has attracted attention. As an epitaxial growth technique of forming a GaN-based semiconductor film, there is an organic metal vapor phase growth method (MOCVD method). In the organic metal vapor phase growth method, for example, organic metal such as trimethyl gallium (TMG), trimethyl indium (TMI), and trimethyl aluminum (TMA) or ammonia (NH3) or the like is used as a source gas.\nJP-A H10-158843 and JP-A 2002-212735 disclose vapor phase growth apparatuses including a plurality of reaction chambers for improving productivity."}
{"text": "1. Field of the Invention\nThis invention relates to a method for the production of alkali metal phosphates and phosphoric acid by the acidulation of phosphate rock and more particularly to a method for the acidulation of phosphate rock in the substantial absence of fluorine evolution, recovery of the fluorides as useful solid values, and recovery of fluoride-free products.\n2. Description of the Prior Art\nPhosphoric acid plants are currently operated utilizing a basic and well known process for the acidulation of phosphate rock which comprises reaction of the rock with sulfuric acid to form phosphoric acid with subsequent reaction of the phosphoric acid, with for example ammonia to produce monoammonium phosphate (MAP) and diammonium phosphate (DAP). The phosphoric acid formed in this process is called wet process phosphoric acid. In this reaction, a by-product is gypsum having the chemical formula CaSO.sub.4.2H.sub.2 O. Essentially all phosphate rock contains some fluorine, normally in the 3.0 to 4.0% range, and the acidulation reaction usually generates gaseous fluorides.\nIn recent years, both air and water pollution laws and regulations have become more stringent and are now being enforced more rigorously. Operating companies have had many pollution problems with fluorine emission into the atmosphere and with the by-product gypsum from these phosphoric acid plants. Thus, an important problem in the operation of these wet process phosphoric acid plants has been in the expensive methods for handling the large amounts of fluorine compounds which are liberated in the gaseous and aqueous effluents from such plants. In some phosphate complexes from 10,000 to 30,000 tons per year of fluorine compounds may be liberated by various methods. It is estimated that in a typical wet process phosphoric acid plant, a portion of the fluorides are evolved into the atmosphere in gaseous form, such as hydrogen fluoride and silicon tetrafluoride, which can destroy vegetation and affect other facilities in close proximity to the plant if they are not scrubbed out and such scrubbing systems are not always effective. A second portion of the fluorine is found in the gypsum dumps and is subject to leaching into groundwater and streams. Still another portion of the fluorine remains with the final products and when such products are used as when adding fertilizers to the soil, they change the soluble and insoluble fluoride levels. It is only in recent years that studies have been made on the effects of fluorides contained in the final product and indications seem clear that they may have a deleterious effect on the long range producing ability of the soil. See for example Kudzin et al, Chem. Ab., 73, 870534 (1970).\nIn these systems, any conventional wet process phosphoric acid technology accomplishes two primary objectives, namely: (1) phosphate rock acidulation, and (2) the growth of readily filterable calcium sulfate crystals either as the dihydrate (gypsum), or as the hemihydrate. Conventional phosphoric acid technology carries out both of these objectives essentially simultaneously which leads to a number of environmental and purification problems almost immediately. The presence of strong sulfuric acid in the acidulation phase releases fluorides as HF, SiF.sub.4, and/or H.sub.2 SiF.sub.6. This poses serious fluoride containment and subsequent recovery problems. Furthermore, unless excess sulfate levels are carefully and closely controlled, minute gypsum crystals can and will blind rock particles and usually result in poor P.sub.2 O.sub.5 recovery. The presence of free H.sub.2 SiF.sub.6 in the acid system leads to severe scaling and excessive maintenance costs even with improved design features to minimize this effect.\nThere is a great deal of literature and patent art related to attempts to remove the fluoride values from fluorine-containing phosphate rock in the operation of a phosphoric acid plant including methods for suppressing the evolution of fluoride values in the operation of a process and/or attempting to scrub the fluorine from effluent gases and waste water. Two such methods are described in U.S. Pat. Nos. 2,954,275 and 2,976,141 to Carothers et al which use sodium or potassium compounds of suppress the fluorides so that they are concentrated in the gypsum cake. These patents indicate that this is achieved by adding a suppressing amount of an alkali metal salt to the acidulation reaction. However, these processes were conducted in the presence of sulfuric acid in the acidulation reactor and the process had incomplete control on fluoride decomposition and evolution during acidulation.\nOther prior art patents have been noted which have also attempted to overcome the problem of fluorine evolution and the reduction of the amount of fluorine contained in final products. A reference of this type is British Pat. No. 735,086 (1955), which discloses the acidulation of phosphate rock by a two step procedure using a strong mineral acid such as nitric acid or hydrochloric acid. According to this patent, an initial low temperature acidulation at 20.degree.-50.degree. C. is carried out with the addition of an alkali, for example, ammonia or lime, as a precipitating agent in a quantity sufficient to precipitate substantially the whole of the fluorine and other impurities but insufficient to precipitate a substantial amount of calcium phosphates. After precipitation and optional removal of this precipitate, indicated as being calcium fluoride by the patentee, the mixture is heated to 75.degree. C. and the acidulation completed by adding a further quantity of alkaline precipitating agent. This patent, however, does not suggest that the use of the strong mineral acid in even the low temperature procedure is going to necessarily provide sufficient reaction to result in at least some fluorine evolution. Moreover, it is doubtful that the precipitate recovered will be in the form of a calcium fluoride. Therefore, while this patentee appreciated the concept of attempting to remove the fluorides prior to the completion of the acidulation reaction, it is clear from a reading of this patent that he did not accomplish the intended purpose.\nIn a similar process, in U.S. Pat. No. 3,431,096 to Hill et al a process is disclosed for reducing evolution of fluorine values in formation of triple superphosphate fertilizer by reaction of phosphate rock and phosphoric acid wherein ammonia or urea is added to suppress the fluorine evolution. However, in this patent, there is no provision for removal of the fluorine values from the product and therefore even if the fluorine evolution is prevented, the fluorine values will be retained in the resulting product and therefore distributed to the soil when it is used as a fertilizer.\nIn a series of patents issuing from the mid 1940's to early 1960's there are disclosed processes for the defluorination of phosphate rock and the production of defluorinated calcium phosphates. In these U.S. Pat. Nos. 2,337,498; 2,442,969; 2,893,834; and 2,997,367, the defluorination reaction is carried out by subjecting a mixture of phosphate rock, phosphoric acid and an alkali metal material to calcination, that is by reaction at temperatures as high as 1000.degree. to 2200.degree. C. Obviously, under these conditions the fluorine is going to be rapidly evolved or if not evolved certainly will remain in the final product, said to be an animal feed. Therefore, these patents did not provide a solution to the problem of fluorine evolution and of fluorine remaining in the final product.\nTwo additional patents of pertinence to processes of this type are U.S. Pat. Nos. 2,567,227 and 2,728,635 to Miller which disclose the acidulation of phosphate rock with phosphoric acid to form monocalcium phosphate, cooling to solidify the monocalcium phosphate and then converting it to dicalcium phosphate by hydrolysis. In the earlier patent, it is indicated that the fluorine in the rock is vaporized in the system, circulates throughout the system and/or leaves the system with the calcium phosphate. The later patent indicates that the process of U.S. Pat. No. 2,567,227 provided a final calcium phosphate product having a fluorine content too high to be of animal feed grade. The solution to this problem in the later patent was the addition of some dilute sulfuric acid in the acidulation step which would of course lead to additional fluorine evolution during the first step.\nThere are also patents known in the art which indicate that it is known to acidulate phosphate rock with phosphoric acid and to then recover solid monocalcium phosphate by cooling of the resulting solution and recovering the monocalcium phosphate. Processes of this type are disclosed for example in U.S. Pat. Nos. 3,494,735 and 3,645,676. In addition, U.S. Pat. Nos. 3,619,136 and 3,792,151 to Case disclose the reaction of phosphate rock with recycle phosphoric acid at temperatures of about 125.degree.-180.degree. F. (52.degree. to 83.degree. C.) to form a solution of monocalcium phosphate, reacting the latter solution with sulfuric acid to produce phosphoric acid and calcium sulfate, precipitating the calcium sulfate, and recycling a portion of the phosphoric acid to the phosphate rock reaction. These patents point out that under the conditions recited, fluorides are not evolved but remain primarily unreacted and may be found with insoluble materials although a portion remains in the phosphoric acid solution product. Thus the products are going to be contaminated with fluorides. It is also known to react phosphate rock or a solubilized form with sulfuric acid and KHSO.sub.4 in combination with other steps and this reaction is described in U.S. Pat. Nos. 3,697,246 and 3,718,253. None of these patents disclose the unique combination of steps and advantages achieved by the process of this invention.\nOther patents in the phosphoric acid and fluoride art suggest methods for use of scrubbing and recycling plants in an effort to contain or convert the fluorides evolved so that as much as possible of the fluorine can be recovered. Nevertheless, in all of these earlier approaches to the problem, provisions are never made for disposing of the approximately one-tenth or more of fluoride contained in the final product nor are there provisions made for converting the fluorides to useful products.\nGerman Offenlegungschrift 2,257,948 discloses the reaction of K.sub.2 SiF.sub.6 with KOH and lime to ultimately form calcium fluoride. However, the K.sub.2 SiF.sub.6 is not obtained from the acidulation of phosphate rock.\nIt is to be appreciated therefore, that the process of the invention provides an improvement over the processes of the references discussed herein and all other references of which applicant is aware. The present invention provides a system which substantially eliminates the problem of the fluorine evolution in the acidulation of fluorine-containing phosphate rock with phosphoric acid, the recovery of useful products from the acidulation reaction and the recovery of the fluorine contained in the rock in a usable form. Therefore, the present invention provides a unique combination of steps and advantages not appreciated heretofore in the prior art as none of these prior references disclose the unique combination of steps and results of this invention."}
{"text": "The present invention relates to microwave radio frequency waveguide feed systems, and more particularly to a high-isolation, broadband, and polarization diverse circular waveguide feed for reception of Direct Broadcast Satellite (DBS) television and Internet satellite downlink services that operate worldwide.\nThe widespread demand for high-quality video, audio, and data communications via satellite has resulted in the need for additional bandwidth and better cross polarization rejection as well as reduced interference from noise or adjacent frequency operation. As a result, satellite broadcast systems are operating over broader and higher frequency ranges and implementing sophisticated methods to reduce interference and improve the intelligibility of communication signals that limit their operating capability. However, the radio frequency apparatus that operate at higher frequencies and with broader bandwidth require considerable design attention and often result in multiple and complicated waveguide feeds in order to account for electric and magnetic field behavior that exists inside the microwave waveguides that propagate their signals.\nAlso, in order to maintain reliable communication, transmit and receive systems must possess polarization compatibility, which is that property of a radiated wave of an antenna that describes the shape and orientation of the electric field vector as a function of time. Polarization compatibility further complicates the waveguide feed design because electromagnetic energy may be transmitted in arbitrary linear, right-hand circular, left-hand circular, or elliptical polarization.\nIt is well known in the art that square waveguides produce mode patterns that allow high efficiency injection or removal of energy for linear polarized electromagnetic waves using probe coupling, which results in orthogonal linear polarizations of high-isolation needed to reduce noise and unwanted adjacent frequency interference. Satellite systems, however, typically operate with circular polarization, which propagates well in circular waveguides, but generates undesirable cross polarization components and poor isolation when using orthogonal probe coupling methods in planar orientation. To minimize cross polarization components that result in a circular waveguide from the two orthogonal polarizations that comprise the circular polarized wave, elaborate conversion methods are employed to transform circular polarized electromagnetic waves. Polarity converters and filters are methods used to condition the circular polarized wave, but have the disadvantage of being difficult to design and possessing high cost and large size.\nThe present invention is a microwave feed assembly of simple, elegant, and scalable design that incorporates the desirable characteristics of broadband operation, polarization diversity, high-isolation between the orthogonal linear polarizations, low insertion losses, small size, and applicability to a broad family of antennas.\nThe present invention relates to a high-isolation, broadband, and polarization diverse circular waveguide feed for microwave frequency antennas. In one aspect of the invention, the waveguide feed supports transmission or reception of any arbitrary linear, right-hand circular, left-hand circular, or elliptical polarized microwave signal while achieving desirable performance over a wide range of frequencies with small size. In another aspect of the invention, the waveguide feed incorporates high cross-polarization rejection of unwanted TE11 Mode components when operating in arbitrary linear mode. In yet another aspect of the invention, the waveguide feed employs high probe-to-probe isolation for rejection of undesired cross-polarization when operating in circular or elliptical polarization mode. A waveguide feed assembly is disclosed, which comprises a combination of symmetrical shaped conical frustrum waveguide and circular waveguide segments together with a novel arrangement of orthogonal and nonplanar electric field probes and radio frequency impedance posts to achieve high-isolation, broad bandwidth, and polarization diversity.\nIt is an object of the present invention to provide a microwave waveguide feed system that can transmit or receive arbitrary linear, right-hand circular, left-hand circular, or elliptically polarized electromagnetic waves.\nIt is another object of the present invention to provide a microwave waveguide feed system that will support operation over a broad range of frequencies.\nIt is yet another object of the present invention to provide a microwave waveguide feed system with cross polarization rejection greater than 20 dB.\nIt is yet another object of the present invention to provide a microwave waveguide feed system with probe-to-probe isolation greater than 30 dB when rejecting undesired linear cross polarization of the two orthogonal linear polarizations that comprise circular or elliptical polarized electromagnetic waves.\nIt is a feature of the present invention to provide a waveguide assembly that is polarization diverse for operation with arbitrary linear, right-hand circular, left-hand circular, or elliptically polarized electromagnetic waves.\nIt is yet another feature of the present invention to provide a compact, reliable, and simple to manufacture waveguide assembly that uses common materials and is suitable for reflector type antennas used to meet minimal radome swept volume applications by reducing the axial length of the waveguide assembly.\nIt is an advantage of the present invention to provide a waveguide assembly that is low cost, rugged, and applicable to a broad family of microwave antennas.\nIt is another advantage of the present invention to provide a microwave waveguide feed that can operate as a stand-alone microwave antenna system.\nIt is yet another advantage of the present invention to provide a waveguide assembly that incorporates design characteristics that are scalable to any frequency of microwave operation.\nThese and other objects, features, and advantages are disclosed in the specification, figures, and claims of the present invention."}
{"text": "The present invention relates to a connector for insulation displacing termination of ribbon cable having conductors on close centerline spacing and particularly a cable with conductors on 25 mil center.\nConnectors for insulation displacing termination of ribbon cable are well known. These generally employ a housing having a mating face, an opposed cable receiving face, and at least two rows of terminal receiving passages extending between said faces. A plurality of stamped and formed metal terminals are received in respective passages, each terminal having major rolled surfaces and sheared edge surfaces, each terminal having a mating portion toward said mating face and a slotted plate toward said cable receiving face.\nThe slotted plate behaves like two parallel beams. For a given force at a given point on a beam, the deflection is a function of the width, thickness, and properties of the material. As contacts become smaller, these dimensions must be balanced carefully to attain a desired force/deflection ratio. A 10% decrease in beam width can necessitate a 37% increase in material thickness to retain the same force/deflection ratio. Simply using a thicker or stronger material may lead to difficulties in stamping and forming a terminal of the desired configuration. Thus, as the centerline spacing of ribbon cable has decreased, connector design problems have arisen.\nA Microminiature D Connector has recently been developed by ITT Cannon Electric for mass terminating ribbon cable with 30 gage wire conductors on 25 mil centers. See Cabourne, Michael K. \"Mass Termination of 25 Mil Planar Cable with Micro D Connectors\", a paper presented at the 17th Annual Connectors and Interconnection Technology Symposium in Anaheim, Calif. The disclosed connector has slotted plate terminals each with a conductor receiving portion extending above the cable receiving face and a base portion in a respective passage. The conductor receiving portion extends to a pair of insulation piercing points and mutually opposed outer sheared edge surfaces which extend from respective points as diverging straight lines toward the base portion. The base portion likewise has mutually opposed outer sheared edge surfaces which diverge as straight lines, which straight lines are colinear with the outer edge surfaces of the conductor receiving portion. This yields a tapered beam design which provides a normal contact force of about 21/2 pounds on a 30 gage stranded wire received in the conductor receiving portion. However, the outer sheared edge surfaces at the level of wire termination are considerably wider than the centerline spacing of the cable, with result that insulation of neighboring adjacent wires is displaced. Thus, while desired contact force is achieved, the possibility of short circuiting between conductors is introduced."}
{"text": "Medical electronic apparatus has been employed for many years for many types of diagnosis and treatment. One field in which there has been considerable growth recently is electrosurgery, in which a suitable generator provides a high frequency, high voltage current which is transmitted to a small surgical electrode having a thin knife-like tip to be applied to a patient. The patient sits or lies on a patient plate and is grounded thereto, with the plate being connected by a further conductor back to the generator. The relatively extremely small area of contact by the electrode with the patient provides an intense current in a highly localized area, producing a cutting action. The current passes through the patient's body to the patient plate where the area of contact is so great that no burning effect occurs.\nFor cutting purposes, the generator is activated to produce a continuous sine wave signal. However, the same instrument may be used to apply to the wound after cutting in order to produce coagulation. For this purpose the generator may be selectively activated to produce a pulsing signal which produces the desired results. Switching means are available for the operator to selectively control an activating means for causing the generator to produce the desired type of current.\nIn the operating room, a surgeon would like to use the electrosurgical device in a fashion similar to that in which he uses a standard sharp scalpel. The surgeon typically grips the scalpel like a pencil with index finger guiding and hear the cutting tip for precise control. Currently available electrosurgical hand instruments, called \"pencils\", are not readily gripped this way, because the mode switches are positioned far from the cutting tip, toward the rear of the instrument. This requires the surgeon to grip the pencil far back on the handle in order to have the index finger positioned to activate the mode buttons, resulting in awkward feel and loss of precise control. An additional disadvantage of electrosurgical pencils is the possibility of inadvertent activations of the mode switches by pressures being applied to the unit. If this occurs when the tip is in contact with the patient or with conductive items, a serious burn can result. Another drawback to the use of electrosurgical pencils has been that many surgeons still desire to use a sharp scalpel for incision-making, with coagulation being achieved by the electrosurgical pencil. This requires the surgeon to continually switch from one instrument to the other.\nElectrosurgical pencils with a rocker switch for selection of mode of operation are disclosed in U.S. Pat. No. 4,228,800 to Degler, Jr., et al, U.S. Pat. No. 3,648,001 to Anderson, et al, U.S. Pat. No. 4,443,935 to Zamba, et al, and U.S. Pat. No. 3,801,766 to Morrison, Jr. As is typical among such devices, the pencil in each is connected to an electrosurgical generating apparatus by a cable comprising three conductors: a first signal line, a second signal line, and a common line which serves not only to return the selected control signal to the generator, but also to deliver current to the scalpel electrode. Thus, all three conductors are maintained at high voltage, with the control signal conductors differing in voltage slightly from the common conductor so that a relatively small current will serve to select the mode of operation desired.\nU.S. Pat. No. 3,720,896 to Beierlein discloses a magnetic slide switch with three positions, and appears to provide electrical isolation when in the center position. U.S. Pat. No. 4,463,759 to Garito, et al. also provides electrical isolation."}
{"text": "Internet electronic mail (hereinafter “email”) is growing at a phenomenal rate. Users often have multiple email accounts with multiple post offices and must log into each account to obtain or check email messages residing at each email service provider. Email clients generally handle this task for users. The process of logging into multiple email servers is generally performed as a background task as long as the user's computer is connected to the Internet. However, a user is generally not notified that new email messages exist unless the user is connected to the Internet. Therefore, the user must frequently connect to the Internet in order to receive new email in a timely manner. The task of checking for new email is generally left up to the client email software.\nSubscribers and users of wireless networks have an even greater challenge. Setting up the equipment to make a connection is time consuming and, therefore, users generally wait until they have email to send before checking for new email. Once connected to the Internet, the client email software checks for any new email.\nAccordingly, there is a need in the art for a system and method for determining when new email is sent to a subscriber, and for sending an email notification to the subscriber alerting the subscriber of the new email. Furthermore, there is a need in the art for sending a wireless email notification to the subscriber and thus providing the subscriber with quicker and timelier access to email. There is a further need in the art to provide wireless email notification to a subscriber and sending that notification to a device that the subscriber uses to connect to the Internet or an internal network and set up that device to automatically make a connection back to the email server."}
{"text": "The present invention relates in general to telecommunication networks, and is particularly directed to a multi-mode backplane architecture, that is configured to support both multiplexed communication circuits and non-multiplexed (loop repeater) circuits in the same equipment shelf, and thereby provide either or both types of telecommunication connectivity between the network and subscriber circuits served by the network.\nPrior to the advent of multiplexedxe2x80x94digital communications, telephony service providers provided telecommunication connectivity to their subscribers by way of dedicated, point-to-point analog copper wire loops, extending from line cards installed in an equipment bank of a service, facility, such as a central office to telephone equipment installed at a customer site. With continuing improvements in digital communication equipment and transport protocols, these conventional point-to-point loops have been augmented by time division multiplexed (TDM) digital circuits, that are capable of transporting a wide variety of communication signals (audio/voice, video and data) to multiple customer premise equipments.\nBecause multiplexed digital circuits require a backplane architecture different than that employed for legacy loop circuit line card equipment racks, telecommunication service providers currently employ two respectively different types of line card equipment shelvesxe2x80x94one for legacy loop circuits; the other for multiplexed circuits. Now even though it is anticipated that the number of multiplexed equipment installations will continue to grow, and may eventually become the predominant mode of network communication circuit, it is still necessary that each telecommunication equipment installation have the capability of providing service for both types of circuits. This means that each line circuit installation facility must be equipped with both multiplexed and loop repeater mode backplane shelves, irrespective of the degree to which each is populated, which increases both the complexity and cost of the installation.\nIn accordance with the present invention, the above problem is effectively obviated by a new and improved multi-mode backplane architecture, that is configured to simultaneously support and provide telecommunication connectivity for both multiplexed communication circuits and non-multiplexed (loop repeater) circuits in the same equipment shelf, so that both signalling modes are available in the same equipment shelf between the network and subscriber circuits connected to the backplane. For this purpose, the backplane architecture of the present invention comprises a set of redundant (for improved reliability) network interface multiplexers, each of which provides both time division multiplexing and point-to-point multiplexing connectivity between the network and line access modules to which subscriber circuits served by the network are selectively ported.\nOn the network side, each redundant network interface multiplexer is ported to a pair of high speed multiplexed network communication links that transport standard high speed multiplexed digital telecommunication signals, such as but not limited to DSX, STS-1, OC-X signals. On the line circuit side, each network interface multiplexer is coupled to a time division multiplexed communication bus, along which a plurality of line access module card slots are distributed, and to a plurality of point-to-point communication links, that are ported to respective ones of the line access module card slots.\nEach network interface multiplexer contains a data buffer that is ported between each of the pair of high speed multiplexed network communication links and a signal path routing control processor. The routing control processor is coupled to a data matrix, that is ported to respective ones of the line access modules, and a time division multiplexed port, that is coupled to the time division multiplexed communication bus. The signal path routing control processor and the data matrix are controlled by a multiplexer control digital signal processor. Once provisioned by the system controller unit, the multiplexer control processor is operative to control the execution of the intended multiplexed telecommunication signalling functionality between the high speed multiplexed network communication links and the line access modules.\nEach access module card slot is configured to receive a line circuit access module that is programmably configurable to provide a selected mode (multiplexed mode or loop repeater mode) of telecommunication connectivity between the network and a subscriber circuit served by the network. For external loop mode connectivity, a plurality of point-to-point communication links are respectively ported between external loop ports of the access module card slots and associated loop ports of a plurality of non-multiplexed interface connectors. Each non-multiplexed interface connector is configured for connection to non-multiplexed communication links of the network and non-multiplexed communication links (e.g., two-wire loops) of subscriber circuits served by the network.\nA system controller unit, such as a desktop computer terminal, is coupled to each of the redundant network multiplexers and the line circuit access modules for allowing network service personnel to program backplane connectivity configuration parameters into the redundant network multiplexers and the line circuit access modules, and thereby enable any line access module to provide either multiplexed or loop repeater mode telecommunication connectivity between the network and subscriber circuit to which that line access module is ported.\nIn order to provide a selected mode of telecommunication connectivity, such as but not limited to ISDN, HDSL, ADSL, T1 office repeater, fractional T1, optical DS2, DSX-1, DDS, DAML, DAML POTS, etc., between the network and a subscriber circuit served by the network, a respective line circuit signal access module contains a processor controlled input/output multiplexer arrangement, that is buffered between each network interface multiplexer and respective sets of non-multiplexed communication links that are ported to the non-multiplexed interface connectors.\nFor this purpose, a multiplexed signal input buffer has a set of inputs ported to the TDM bus and point-to-point links from the network interface multiplexers, and a set of outputs ported to multiplexed inputs of an input multiplexer portion of the input/output multiplexer arrangement. On the non-multiplexed signalling side, a non-multiplexed signal input buffer has a set of non-multiplexed signal-associated inputs ported to point-to-point loops of the non-multiplexed interface connectors. On its output side, the non-multiplexed signal input buffer has a set of outputs ported to associated input ports of the input multiplexer portion of the input/output multiplexer arrangement. The input multiplexer has a plurality of outputs associated with respective ones of its multiplexed and non-multiplexed inputs coupled to corresponding inputs of a protocol conversion, network/subscriber circuit connectivity-defining, signal processing circuit.\nThe signal processing circuit is preferably implemented as a digital signal processor that is programmably configurable by the system controller unit to execute protocol conversion functionality for effecting one of loop repeater and time division multiplex connectivity between selected ones of its input and output ports, so as to provide a selected mode and type of telecommunication connectivity, such as those referenced above. The signal processing circuit has a plurality of output ports coupled to corresponding inputs of an output multiplexer portion of the input/output multiplexer.\nOn the multiplexed signal side, the line access module\"\"s output multiplexer has TDM output bus ports respectively coupled to TDM bus ports, and point-to-point output bus ports respectively coupled to associated inputs of the multiplexed signal input buffer. Likewise, on the non-multiplexed signal side, the output multiplexer has a set of non-multiplexed output ports coupled to respective non-multiplexed signal-associated input ports of the non-multiplexed signal input buffer. The mode of connectivity through and operation of the line access module\"\"s input/output multiplexer arrangement and its protocol conversion, network/subscriber circuit connectivity-defining, signal processing circuit is controlled by an access module control processor that is coupled to the backplane\"\"s system controller unit. Since each line access module is selectively connectable to and individually programmable for providing connectivity and protocol conversion for either multiplexed or loop mode signalling paths of both network and subscriber connections, the backplane enjoys complete flexibility in terms of types of circuit and mode of operation that may be emulated by each line access module."}
{"text": "Cellular communications systems continue to grow in popularity and have become an integral part of both personal and business communications. Cellular phones allow users to place and receive voice calls most anywhere they travel. Cellular phones and other handheld wireless communication devices typically include a radio, e.g. having a wireless transceiver and associated circuitry connected thereto, and an antenna connected to the radio.\nAntenna noise temperature has been discussed in many books and papers, such as John D. Kraus and Ronald J. Marhefka, “Antennas: For all Applications”, McGraw Hill, 2002, ch. 12; Constantine A. Balanis, “Antenna Theory: Analysis and Design” John Wiley & Sons Inc. 1997, ch. 2; David M. Pozar, “Microwave Engineering”, Addison-Wesley Publishing Company, 1993, ch. 12; J. Dijk, MJeuken and E. J. Maanders, “Antenna noise temperature”, Proc. IEEE, Vol. 115, No. 10, October 1968, pp 1403-1409; and Warren L. Flock and Ernest K. Smith, “Natural Radio Noise—a Mini-Review”, IEEE Trans. on AP Vol. Ap-32, No. 7, July 1984 pp 762-767.\nThe definitions for antenna noise temperature are mainly given based on remote sensing and satellite receiving applications, where antennas are generally physically away or well shielded from radio receivers and high gain antennas are used to capture weak signals. In this case the total noise at the terminal of the receiver antenna is mainly contributed from thermal noise and background noise. In contrast, a wireless handheld antenna is physically very close to its receiver so that the printed circuit board and accessories operate as a part of the antenna. This makes the noise contributions to the handheld wireless device antenna different from the noise contributions to antennas for remote sensing and satellite receiving applications.\nThis difference makes the standard antenna temperature definition inadequate for explaining the receiver behavior of handheld wireless devices in a noisy environment. A wireless handheld device generally operates in an ever changing noise environment, and the handheld antenna radiation pattern is generally a broad beam antenna pattern. Furthermore, human physical interface and device usage scenarios change constantly in the practical application. For these reasons, antenna noise temperature is constantly changing in the practical sense."}
{"text": "The invention relates to a shooting head for a core shooter for shooting foundry cores or molds, the shooting head having an opening at its inlet end for selective connection to a filling device immersing with a discharge end or outlet member into the shooting head and predetermining the filling level, or to a clamping head serving to sealingly apply compressed air, and a shooting plate arranged at its outlet end and comprising shooting nozzles.\nCore shooters with shooting heads of the kind under discussion have been known from practice for a long time. For example, reference may be made to DE 41 28 952 C1.\nHowever, within the scope of an automatic production of core packs, the known shooting heads of the described kind are problematic. Depending on the core being shot, differences in shooting weights from 60 kg to 1 kg occur, the shooting heads being normally capable of receiving core sand in a quantity of at most 70 kg. Regardless of the size of the core being shot, the outside dimensions of the shooting heads are however predetermined, in particular in the automatic production of core packs. This means that the shooting heads in the core shooter must be exchangeable regardless of the cores being shot. Consequently, in the case of shooting very small cores, there is a large supply of sand and binding agents in the interior of the shooting head, the small shooting weight causing only a small throughput of core sand and binding agents. In the presence of large supplies of sand in the shooting head, and in particular as a result of the temperature inside the shooting head, which is about 35.degree. C., this small throughput leads to an ageing of the binding agent that is mixed with the core sand, so that cores shot with a very small throughput of sand often lack adequate strength.\nFurthermore, it is necessary that the compressed air required for the shooting penetrate through the entire sand supply within the shooting head, so that a rather considerable compression occurs already in the shooting head, which complicates the shooting of the cores, or makes it even impossible when lumps form in the region upstream of the nozzles.\nIt has already been proposed in DE 41 28 952 C1 to alleviate the problem of penetrating the sand supply with compressed air inside the shooting head to the extent that a filling device or an outlet member thereof immersing into the shooting head fills the shooting head to a certain, predetermined filling level, thereby filling the shooting head only in part. Furthermore, the filled-in molding sand is flattened more or less by a closing flap of the outlet member, so that the core sand being within the shooting head is quasi evenly distributed. Nonetheless, also in this instance it is necessary to fill in a considerable quantity of core sand, even for shooting very small cores, since the filling with only a small amount of core sand would not supply all shooting nozzles with an adequate amount of sand. A complete filling of the shooting head still leaves, even after a flattening, the aforesaid ageing problem on the one hand and, on the other, the problem of penetrating the total supply of core sand in the shooting head with compressed air, so that the risk of unwanted preliminary compressions within the shooting head is incurred likewise in this instance.\nIn its not-yet published German Application P 42 08 647.7-24, applicant has already proposed itself a further solution to the problem with respect to the risk of a preliminary compression occurring inside the shooting head. To avoid an unwanted preliminary compression within the shooting head, a device is provided for generating a turbulence or for an even distribution of the compressed air flowing from the air supply into the shooting head, the device being a so-called streamline filter with slot-shaped passages. However, the ageing problem is thereby not solved.\nEven a synopsis of the prior art known, on the one hand, from DE 41 28 952 C1 and relating to a filling device, and the turbulence of the compressed air required for the shooting by means of a streamline filter as claimed on the other hand in German Application P 42 08 647.7-24, will be able to alleviate at most, but will not solve the problem that arises in particular with the shooting of very small cores both with respect to the ageing of the binding agent within the shooting head and with respect to an unwanted preliminary compression within the shooting head.\nIt is therefore the object of the present invention to describe a shooting head for a core shooter for shooting foundry cores or molds, which is also suitable in particular for shooting very small cores."}
{"text": "Switching structures or switching networks operating according to an asynchronous transfer mode are provided for future communication systems, particularly broadband networks. In this asynchronous transfer mode, packets having a fixed length, these being referred to as cells in the technical field, are transmitted in uninterrupted fashion on every transmission section. These cells are composed of 48 octets for the useful information and of 5 octets of cell routing information for the cell header. When no useful information is to be transmitted at the moment, specifically marked idle signals are transmitted.\nA switching network operating according to the asynchronous transfer mode is structured into one or more switching stages, similar to the switching network of traditional digital communication systems, whereby the individual switching stages are meshed with one another such that a largely blocking-free switching structure arises. The input and output lines of such a switching structure are identified from the cell routing information recited in the cell header of the cell, i.e. from virtual routing information and virtual channel information. These switching-oriented information are derived from the multi-addressing and destination information signalized by the subscribers of a communication system, i.e. telephone numbers of the destination subscribers.\nIt is known from the publication \"ATM-Technologie fuer zukuenftige Breit-bandnetze\" of Siemens AG, page 17, right-hand column, to additionally identify routing information for the switching of a cell through a switching structure from the cell routing information recited in the cell header of the cell and to attach this to the cell. With the assistance of these internal routing information, the respective cell seeks the defined path through an n-stage switching structure under self-control, i.e. without the cooperation of a central communication system controller. As used herein the term \"self-control\" refers to the self-routing of cells. The additional, internal routing information are inserted into an internal cell header of the communication system and are set in front of the respective cell. After the self-switching of the cell through the switching structure, the internal cell header of the communication system is again removed.\nThe data scope of the cell routing information or of the internal cell header of the communication system is essentially defined by the plurality of switching stages, whereby one information string, for example, one byte is usually provided for a switching stage, these being compiled to form a block. Internal cell headers of the communication system are being defined to an increasing degree that have a prescribed plurality of information strings matched to the maximum plurality of switching stages, for example 5 bit octets or bytes for 5 switching stages.\nIn the case of multi-address transmission of cells to a plurality of output lines of the switching structure, the self-control principle of the cells can no longer be maintained without further measures given internal cell headers of the communication system structured in this way."}
{"text": "In commercial kitchens, the steam and vapor carrying latent heat, and the sensible heat produced within the kitchen must be carried off by appropriate ventilating means to maintain reasonable temperature and humidity conditions within the kitchen so that these conditions are tolerable by the kitchen personnel and do not affect adversely the cooking process.\nIn other words, the ventilation must be so dimensioned that predetermined temperature and moisture levels in the kitchen are not exceeded.\nThis requires generally a certain air replacement rate with which fresh outside air is supplied to the kitchen and kitchen air with the sensible heat and steam and vapor with the latent heat are withdrawn.\nDuring periods in which the facility must be heated, generally in the winter, the supplied air must be heated to prevent the kitchen personnel from being chilled and to prevent detrimental temperature drops which might interfere with the cooking process. Furthermore, because relatively large amounts of air must be circulated, the heating may be relatively costly with respect to the energy which must be supplied to the air.\nIn determining the amount of air which must be circulated, therefore, one must consider on the one hand the cooking units which are in operation in the kitchen and may contribute comparatively large amounts of water vapor and hence latent heat, namely, boiling equipment, frying and water baths and steam tables which are commonly present in such kitchens. It is also necessary to consider cooking devices which contribute primarily radiant heat and/or convection heat to the atmosphere in the kitchen, i.e. sensible heat. These can include grills, ovens, tilting pans and the like. The type of heat which makes the most significant contributions to the air circulation, therefore, will depend on the facilities in the kitchen.\nTo minimize the amount of heat energy which must be supplied for the heating of the cold supplied fresh air, two basically different approaches are employed. In one approach heat is removed from the discharged air and used to preheat the incoming air by heat exchangers. A drawback of this system is that the discharged air normally contains residual grease which can condense in the heat exchangers and requires maintenance thereof at high cost. The other approach utilizes a so-called short circuit or induction hood in which a part of the supplied air is blown via a nozzle directly into the vapor collection space and thus is immediately drawn out in a short circuiting of the fresh air supply. In winter in which less sensible heat must be extracted, this short circuiting allows less flow from the kitchen of already heated air while nevertheless supplying a sufficient velocity with the aid of the unheated external air to permit the grease filter in the hood to remain effective and maintain the high degree of grease removal which is necessary. It should be understood, in this regard that the grease filter generally operates as an impingement or baffle filter in which the relatively high velocity is necessary for a high degree of grease separation.\nThe drawback of this approach, however, is that the introduction of unheated outside air into the kitchen and hood requires expensive and careful control of the air intake because any of this air which is not immediately discharged will alter the climate within the kitchen."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for sensing a human body that is available, for instance, in an automatic door or the like.\n2. Description of the Prior Art\nIn an automatic door, approach of a human body to the door is sensed, and in response to the sensing of a human body a door opening/closing signal is generated to actuate the door to open and close, and to that end, as a method and an apparatus for sensing a human body, various kinds of methods and apparatus have been known in the prior art.\nFor example, a method and an apparatus for sensing a human body, in which a projector for projecting infra-red rays and a photo-sensor adapted to receive infra-red rays reflected by a floor surface and a human body are provided and the photosensor generates a sense signal when it senses variation of the amount of reflection of infra-red rays, have been known.\nHowever, these heretofore known methods and apparatus had the following shortcomings:\n(1) If rays of sunlight should momentarily enter the photo-sensor as reflected by any moving body, then a sense signal would be incorrectly generated, and hence sometimes, a door may malfunction. PA1 (2) If the infra-red rays projected from the projector should momentarily enter the photo-sensor as reflected by falling snow, then a sense signal would be incorrectly generated, and hence sometimes, a door may malfunction. PA1 (3) Under the condition where a human body stands still, it cannot be sensed.\nMore particularly, since the amount of reflection from a floor surface of the irradiated infra-red rays would change according to variation in time of a radiation efficiency of the projector and according to variation of the floor surface condition, in the event that a level of the amount of reflection of the infra-red rays is simply used for determination of sensing of a human body, minute variation in time of the level itself of the amount of reflection would be caused by the above-mentioned variation of the amount of reflection, hence it is necessary to inhibit sensing of that variation, consequently a level difference of a minute amount of reflection becomes hard to be senced, and so a sensing distance cannot be chosen long.\nIn order to resolve this problem, a method in which the amount of variation of the reflection amount is sensed rather than the reflection amount itself and a sense signal is provided depending upon the amount of variation, that is, a method of differential operation type, may be employed, but if such type of method is employed, in the case where a human body stands still, the human body cannot be sensed because the reflection amount does not vary."}
{"text": "As illustrated in FIG. 1, SRAMs (static random access memories) are characterized by having gate patterns 15 aligned, in a broken-line manner, perpendicularly to active regions 18. According to the example of FIG. 1, the SRAM has point-symmetric cell structures, and two transfer transistors and two CMOS (complementary metal oxide semiconductors) inverters are provided in each cell 100 symmetrically around a point.\nIn reducing the size of a SRAM, it becomes a key issue of how much a protruding amount B of each gate pattern 15 from the active region 18 can be reduced, as illustrated in the dashed-line box A. Next is described a current problem with a focus on, for example, driver transistors of the SRAM of FIG. 1.\nFIG. 2 is an enlarged view of the region A of FIG. 1, and illustrates the setback position of a gate end portion created in a gate etching process. In general, end portions of gates 25 formed by the gate etching process are located in setback positions in resist patterns (gate patterns) 15. Therefore, the gate protruding amount B needs to be sufficiently provided in advance when the resist patterns are formed, in view of the setback amount of the gate etching. This, in turn, requires providing sufficient spacing “d” between the active regions 18 in view of the setback amount of the gate etching, which prevents a reduction in the size of the SRAM device.\nFIG. 3 illustrates a setback of the gate end portion after the gate etching and device failure. In the case where the gate protruding amount B (see FIG. 1) is sufficiently provided, the source and the drain are separated by the gate, as illustrated in FIG. 3A, and therefore, a favorable transistor may be formed. However, if the gate protruding amount is insufficient, the gate end is positioned posteriorly by exposure of polysilicon during the patterning process and the gate etching. As a result, as illustrated in FIGS. 3B and 3C, the gate end portion does not sufficiently overlap the active region (the source and drain). In particular, in the case of FIG. 3C, the source and the drain are not separated by the gate, causing a short circuit, and thus, the device is completely defective. In the case of FIG. 3B, although the source and the drain are separated by the gate and the sidewall, the gate length is different from that of the favorable device (FIG. 3A). Accordingly, there are differences in the device properties, and therefore, the device of FIG. 3B is also determined as defective.\nThe above description is given with an example of a driver transistor near a cell boundary; however, the same problem may occur for the protruding amounts of the transfer gates within the cells of FIG. 1.\nGate double patterning has recently attracted attention as a technology for preventing setbacks of gate etching end portions and decreasing cell sizes of SRAMs by reducing the space “d” between the active regions 18 of FIG. 2 (For example, see M. Kanda, et al, “Highly Stable 65 nm Node (CMOS5) 0.56 μm2 SRAM Cell Design for Very Low Operation Voltage”, 2003 Symposium on VLSI Technology Digest of Technical Papers, pp. 13-14). According to the technology, a single long gate pattern connecting adjacent gates is created first, and then etching is performed using a gate-separating mask 20 having an aperture 21 so as to form separated gates, as illustrated in FIGS. 4A through 4C. The technology does not cause setbacks of the gate end portions, and therefore, it is possible to reduce the space “d” between the active regions 18 of FIG. 2.\nHowever, the inventors of the present disclosure have found a problem associated with the gate double patterning of FIGS. 4A through 4C. If the gate 25 is cut at a position very close to the active region (the source and drain region) due to displacement of the gate-separating mask 20 during the exposure process, as illustrated in FIG. 5A, and a device is then created according to general procedures, changes occur in the current characteristics of the gate end portion.\nFor example, four-way angled implantation is performed in order to form a pocket 26, as illustrated in FIG. 5B, then a sidewall (SW) 27 is formed by extension implantation, and a source-drain region 28 is formed, as illustrated in FIG. 5C. In this case, ion implantation characteristics are different between a region adjacent to the edge along the gate end portion and the remaining region. Therefore, variation is caused in the current characteristics (the arrow b) close to the edge of the gate and the current characteristics (the arrow a) of the inside the gate.\nDevices produced in this manner cause property fluctuations, which become a factor of being defective. In order to prevent such property fluctuations, it is necessary to provide a sufficient amount of spacing “d” between the active regions 18 of FIG. 2 in view of a margin of exposure displacement of the gate-separating mask and an implantation margin even in the case of performing the gate double patterning."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of controlled flow, exhaust manifold systems and, more particularly, to apparatus for controlling the flow of exhaust gases between a plurality of rocket stations and a common exhaust gas manifold or plenum tube connected thereto.\n2. Description of the Prior Art\nIn many military applications, numbers of rockets are stored or disposed in closely adjacent magazine chambers, launch tubes, or the like, hereinafter referred to collectively as chambers. Exhaust gas outlets are normally provided, even from magazine storage chambers, to duct rocket exhaust gases generated during intended or accidental rocket ignitions to a safe location. Where available space is at a premium, for example on ships, manifolding of a number of chambers into a common exhaust duct or plenum tube is often necessary.\nObvious problems exist if ducts connecting the chambers to the common exhaust manifold are always or normally open. When one (or more) of the rockets is intentionally or accidentally ignited, at least portions of the resulting exhaust gases, which may be at about 6000.degree. F., will be circulated through the common manifold and into other chambers through the open connecting ducts. Rockets and rocket warheads in these other chambers are very likely to be ignited or be detonated by these hot exhaust gases. If these other rocket chambers are open at upper ends, as are launch tubes and some storage compartments, exhaust gases entering the chambers through the connecting ducts escape through the open ends and may cause extensive heat damage to adjacent installations.\nTo prevent such occurrences, some type of safety door or gas valve is normally installed either at the outlet opening of each rocket chamber or in the connecting duct to the exhaust manifold. When a rocket is accidentally or intentionally ignited, the associated safety door or gas valve is caused to open--usually by the exhaust blast--to admit the exhaust gases into the manifold. The doors or valves associated with other chambers are maintained in a closed condition to prevent circulation of the exhaust gases thereinto.\nThe patent disclosures, for example, of Eastman and Neuman et al. (U.S. Pat. Nos. 2,445,423 and 3,228,286, respectively) illustrate use of such doors or valves. Previously available or disclosed apparatus, however, have substantial disadvantages. For example, the patent of Neuman et al discloses at the bottom of each compartment of a multiple rocket storage magazine, a non-hinged, \"blow out\" door. These doors lead through conducting ducts to a common exhaust manifold. If any of the rockets in the magazine are accidentally ignited (for example by enemy fire), the force exerted by the resulting rocket exhaust gas on the upper surface of an associated door blasts the door out of its opening and admits the gases into the manifold. An associated fire extinguishing system is designed to direct pressurized water through the resulting opening and extinguish the rocket. A major disadvantage, however, is that no means are provided for automatically reclosing the door after the rocket has been extinguished. Unless the blow-out door is manually replaced--for which little provision seems to have been made--hot exhaust gases from subsequent accidental ignition of another rocket would enter the compartment and could cause reignition of the rocket or explosion of its warhead before such next-firing rocket is extinguished. In addition, if the compartments are not sealed in upper regions--which they do not seem to be--hot exhaust gases from the next firing rocket would be conducted through any compartments containing previously ignited rockets and directly to the rocket launching platform positioned just above the magazine.\nAnother very substantial problem associated with the apparatus disclosed by Neuman et al, and other similar apparatus, is that little consideration appears to have been given to preventing recirculation of exhaust gases back into and through a chamber while a rocket is firing in that chamber. Whatever type of exhaust flow control door or valve is used, it must be suitably configured to prevent exhaust gases emitted therethrough and into the exhaust manifold from flowing around the exhaust stream and back into the rocket compartment. If this occurs, the gases may cause structural damage to portions of the rocket, ignition of other propellants (if the rocket has other stages) or detonation of the rocket warhead. Ignition of these other propellants or detonation of the warhead could ignite or detonate adjacent rockets and warheads, thereby initiating a disastrous chain reaction.\nMerely to provide properly opening and closing rocket exhaust gas flow control doors is, therefore, insufficient; the doors must be configured so that at all exhaust flow conditions they will open only that amount which will cause the rocket exhaust stream to function as a complete \"gas plug\" in the opening to prevent recirculation of exhaust gases back into the chamber.\nThe patent of Eastman discloses apparatus adapted for storing a number of rockets, wherein exhaust nozzles of the rockets are seated in sealing relationship upon short ducts or nozzle extensions leading to a common exhaust manifold. Toggle clamps are used to hold the noses of the rockets in the storage apparatus and no actual storage compartments are formed. Each nozzle extension has, at its lower end, a pair of hinged doors, spring biased to a normally closed condition. Exhaust gas pressure from an accidentally ignited rocket forces the associated nozzle extension doors to swing open against the springs, thereby admitting the gases to the manifold, from which they are discharged at a remote location. The resulting gas pressure in the manifold acts upon under sides of other closed doors to force them tightly closed and prevent circulation of hot exhaust gases into the other nozzle extensions.\nHowever, the door hinges and biasing springs are positioned directly in the path of hot exhaust gas flow from an above firing rocket and will receive maximum heating and erosion therefrom. As a result of heat and erosion damage, the doors immediately below a firing rocket, even if not burned completely loose, as is likely, would probably fail to return to the closed condition after the firing. Also, very possibly, heat from hot exhaust gases flowing through the manifold would damage the biasing springs of other doors. Even if these doors were kept closed by pressure in the manifold during that particular firing, they might subsequently sag open. Then, upon a next accidental rocket firing, the flow of gases through the manifold could force the sagging doors open, rather than closed, allowing circulation of the hot gases into above nozzle extensions with consequent ignition of the associated rockets.\nEven though spring-loaded flow control doors might be satisfactory for use associated with storage of small rockets, wherein firing is unlikely and when it occurs the firing time is short, such doors would be entirely unsatisfactory in applications in which they would be subjected to repeated or sustained rocket exhaust gas flows. They would thus be unsatisfactory for use associated with storing or launching large rockets or with launch tubes from which a large number of even small rockets would be fired.\nFor these and other reasons, improvements in controlling flow of rocket exhaust gases associated with a plurality of rocket stations and a common exhaust gas manifold are not only desirable, but necessary."}
{"text": "The semiconductor devices that utilize the effect of carrier recombination to produce the electromagnetic radiation are widely used in electronics and optoelectronics. Typically, to achieve efficient carrier recombination, the portions of semiconductor materials are brought in contact that provide two different general conductivity types, known as electron and hole conductivity types. Due to the mutual carrier penetration through the said contact caused by diffusion and/or drift in the externally applied electric field, the two types of carriers are simultaneously present in the region of the semiconductor device adjacent to the said contact, also referred to as active region, where they can recombine with each other under the emission of energy, in a form of either electromagnetic radiation (photon) or elastic crystalline lattice vibration (phonon). For a large class of devices, the said photon generation is the useful output.\nIn most of the cases, the photons are produced most effectively if the electron and hole interact and recombine as “free” carriers, resulting in so-called band-to-band recombination. This recombination leads to the generation of a photon with energy of (slightly higher than) the electronic bandgap of the portion of semiconductor comprising the said active region of the semiconductor device. This photon can be effectively absorbed by the semiconductor material having the same or lower bandgap as/than the said active region of the device. To prevent this absorption, efficient photon generating semiconductor devices are usually designed in such a way that the active region comprises the semiconductor material of the lowest bandgap all over the said device. The semiconductor device comprising the regions of different semiconductor materials is usually referred to as a heterostructure. Thus, effective photon generation device is usually a heterostructure device.\nExactly one electron and one hole contribute to a single act of recombination with generation of a photon. The probability of the recombination is proportional to both electron and hole concentrations at a given position within the active region of the semiconductor device. The steady-state concentration of the carriers during the device operation is achieved by detailed equilibrium of the electron and hole concentrations vanishing due to recombination and the electron and hole electric currents into and from the said active region.\nThe narrower bandgap semiconductor material usually has both conduction band and valence band within the energy range of the bandgap of the surrounding higher bandgap semiconductor material. Thus, the active region of a properly designed photon generating semiconductor device normally comprises one or several potential wells for both types of carriers. This helps maintain the carrier confinement within the said active region, increasing the photon generation efficiency of the device. However, the ability of the narrow band gap material comprising the active region to confine the carriers becomes weak with increasing carrier concentrations. This happens mainly because higher electric currents are needed to produce higher steady state carrier concentrations; these higher currents require higher electric fields in and near the active regions; carriers receive extra energy from the higher electric fields and obtain higher probability to drift into higher bandgap regions, thus leaving or passing through the active region. Typically, due to lower effective mass, the electrons gain more energy from the electric fields and are more likely to leave the active region. Due to current continuity, the number of the electrons and holes entering the active region in unit time is equal; as a result of the electrons having higher probability to escape, the symmetry of carrier concentration in the active region vanishes, and the efficiency of photon generation degrades. In addition, at high injection levels, the effect of a limited rate of radiative recombination through any centers involved in the desired radiative recombination may become relevant. This for example includes the considerations of a dipole matrix element reduced by the presence of an electric field within the quantum wells. It also may include the finite rate of excess carrier relaxation and thermalization processes required for carriers to reach the centers of desired radiative recombination.\nIn order to prevent the carrier escape from the active region, additional barrier regions are sometimes introduced adjacent to the said active region. Those barrier regions are usually designed to block the electron escape more effectively and, therefore, are usually positioned within the portion of the said semiconductor device having hole conductivity type. The pay-off of the barrier regions introduction is usually the degradation of the hole injection into the active region and increase in the electric potential drop across the device necessary to achieve the desired level of photon generation. In addition, as a result of the hole injection degradation, the active region becomes flooded with the electrons and, therefore, is turned into the electron (n-) conductivity type, at least for the conditions corresponding to the photon generation.\nIn spite of the barrier regions introduction, the carriers' leakage out of the active region still becomes severe at high carrier concentrations in the said active region. It happens mainly because the external electric potential applied to the device drops mostly across the less conductive regions of the device, which are in the case described above the barrier regions. At high enough applied electric potentials, the barriers become punched through for the electrons escape. This limits the efficiency of the photon generation by the semiconductor device at high currents and even leads to the saturation of the total number of photons that the said device is able to produce.\nThe present invention solves the problems of the electrons leakage from the active region of the photon generating semiconductor device and of balancing the electron and hole concentrations in the active region(s) of the said device by taking independent control of the electric potential applied between the active region of the device and the regions of the said device of different conductivity types."}
{"text": "The invention relates to apparatus and methods for performing automated assembly operations such as riveting, tacking, or the like, on large panel-shaped workpieces, using positioners that manipulate tool-supporting yokes or frames on which assembly tools are supported.\nLarge panel-shaped structures used in the construction of airframes commonly have reinforcing members, usually called stringers, attached to outer skins so as to form strong yet lightweight panels. The stringers are usually affixed to the skins by riveting. In a commonly used assembly technique, opposed riveting tools are positioned by two spaced legs of a yoke or C-frame such that the tools are on opposite sides of a workpiece. The C-frame is carried by a multi-axis positioner, which moves the C-frame along the workpiece for riveting at various locations of the workpiece. Typically, the C-frame can be positioned adjacent only one edge of the workpiece, and the legs of the C-frame are used to reach widthwise across the workpiece. In this type of assembly apparatus, if riveting is to be done near the opposite edge of the workpiece, the throat of the C-frame must have a depth about equal to the maximum width of the widest workpiece that is to be processed in the apparatus. Consequently, the C-frame becomes large and heavy. Because the C-frame has a high mass and inertia, the foundation and traversing mechanism that support and move the positioner on which the C-frame is carried must be made very substantial. Accordingly, the overall cost of the assembly apparatus is relatively high. The high inertia of the C-frame and positioner also means that the maximum acceleration of the positioner is relatively low, and thus the apparatus is slow.\nAs an alternative to sizing the C-frame large enough to reach across the full width of the workpiece, the C-frame positioner can be movable so that it can be positioned at either edge of the workpiece, or the workpiece can be turned 180xc2x0 so that each edge of the workpiece can be positioned adjacent the C-frame. The former possibility necessitates a significantly more complicated positioning system, and the latter possibility is cumbersome and inefficient, particularly where very large workpieces, such as wing panels or the like, are being processed.\nA further drawback associated with C-frame type assembly machines, and other types of assembly machines using positioners for manipulating tool-supporting frames (whether C-shaped or of other shapes) that support assembly tooling, is that when the tool-supporting frame and/or its associated assembly tooling require maintenance or tool changes, the frame and its positioner must generally be positioned in a location near the workpiece, and this prevents other positioners from accessing the area of the workpiece near the inactive frame to continue the assembly operations. This same problem also affects gantry-type machines and other types of frame-based machines. Thus, while maintenance or tool changes are being performed, the utilization rate for the assembly cell drops significantly.\nThe above needs are met and other advantages are achieved by the present invention, which provides assembly apparatus and methods that, in a first aspect of the invention, enable the entire area of a panel-shaped workpiece to be accessed by assembly tooling while avoiding the necessity of using large and heavy C-frames.\nTo this end, an assembly apparatus in accordance with a preferred embodiment of the invention includes a bed defining a support surface extending along orthogonal X- and Y-axes, and workpiece supports adapted to support a workpiece in a position spaced from the support surface with opposite first and second longitudinal edges of the workpiece extending generally parallel to the X-axis and a width of the workpiece extending generally parallel to the Y-axis. The apparatus also includes at least first and second positioners each supporting and positioning a C-frame having a pair of spaced legs adapted to support tooling for operating upon the workpiece when the workpiece is positioned between the legs. Each positioner comprises a platform operable to travel along the support surface of the bed in a direction parallel to the X-axis and a carriage on which the respective C-frame is mounted and which is supported on the platform so as to be movable relative thereto. The first positioner is arranged to carry its respective C-frame adjacent the first longitudinal edge of the workpiece such that free ends of the legs extend toward the opposite second longitudinal edge of the workpiece. The second positioner is arranged to carry its respective C-frame adjacent the second longitudinal edge of the workpiece such that free ends of the legs extend toward the first longitudinal edge of the workpiece. The legs of each C-frame have sufficient length to carry the tooling across at least about half of a maximum width of the workpiece, and the carriages are movable generally away from each other on the respective platforms to enable the positioners to pass by each other as they travel along the bed parallel to the X-axis. Thus, each positioner can access the entire length of the workpiece with interference from the opposite positioner. Because each C-frame need be no longer than necessary to reach halfway across the workpiece, each positioner and C-frame can be made relatively lightweight and, therefore, can be moved relatively quickly from one position to another.\nIn another aspect of the invention, apparatus and methods are provided enabling improved utilization rate of an assembly cell by allowing maintenance and/or tooling changes to be made to a tool-supporting frame in a location remote from the associated positioner, so that the positioner can be used with a replacement tool-supporting frame while the original tool-supporting frame is being serviced. To these ends, an apparatus in accordance with a preferred embodiment of the invention comprises a bed defining a support surface extending along orthogonal X- and Y-axes, and workpiece supports adapted to support a workpiece in a position spaced from the support surface with opposite first and second longitudinal edges of the workpiece extending generally parallel to the X-axis and a width of the workpiece extending generally parallel to the Y-axis. The apparatus also includes a first positioner supporting and positioning a tool-supporting frame that supports tooling for operating upon the workpiece. The first positioner comprises a platform operable to travel along the support surface of the bed in a direction parallel to the X-axis and a carriage on which the tool-supporting frame is mounted and which is supported on the platform so as to be movable relative thereto. The tool-supporting frame is releasably mounted on the carriage such that the tool-supporting frame can be removed and replaced. The apparatus further includes at least one additional replacement tool-supporting frame, a service station proximate the bed for storing tool-supporting frames and for performing work operations thereon, and a frame changer operable to remove the tool-supporting frame from the carriage, transport the replacement tool-supporting frame from the service station to the carriage and install the replacement tool-supporting frame into the carriage, and transport the removed tool-supporting frame to the service station such that work can be performed thereon.\nThe frame changer advantageously comprises a ground vehicle operable to travel to the service station and pick up the replacement tool-supporting frame, travel to a position proximate the positioner having the tool-supporting frame to be replaced, remove the tool-supporting frame from the carriage of the positioner, install the replacement tool-supporting frame into the carriage, travel back to the service station, and unload the replaced tool-supporting frame at the service station. The vehicle can be guided in its movements by any suitable position sensors and a controller.\nPreferably, the assembly apparatus includes at least a second positioner supporting a carriage and a tool-supporting frame releasably retained in the carriage, whereby the second positioner can continue to perform work operations on the workpiece while the frame changer is removing and replacing the tool-supporting frame on the first positioner, and vice versa. The first and second positioners are arranged to carry the respective tool-supporting frames along opposite sides of the workpiece, and are preferably configured so they can pass each other as they travel in the lengthwise direction of the workpiece. Depending on the size of the workpiece, it may be advantageous to employ more than one positioner on a given side of the workpiece. Regardless of the number of positioners used, the invention enables improved efficiency and utilization rate by allowing service to be performed on the tool-supporting frames in an xe2x80x9coff-linexe2x80x9d fashion so that positioners are not tied up while the servicing is being performed.\nThe invention also provides a method for performing assembly operations on panel-shaped workpieces. A preferred method of the invention comprises engaging the workpiece with tooling members supported on a first tool-supporting frame so as to perform assembly operations on the workpiece; moving the first tool-supporting frame relative to the workpiece by a positioner, the positioner releasably supporting the first tool-supporting frame and being operable to travel along the bed to move the first tool-supporting frame lengthwise along the workpiece such that assembly operations are performed at a plurality of locations on the workpiece; taking the first tool-supporting frame out of action so that service can be performed thereon, by moving a frame changer into a position adjacent the positioner, activating a transfer mechanism to cause the first tool-supporting frame to be released from the positioner and to be transferred onto the frame changer and to cause a second tool-supporting frame to be transferred from the frame changer and to be installed into the positioner; moving the frame changer from the position adjacent the positioner to a service station remote therefrom, and transferring the first tool-supporting frame from the frame changer to the service station; and performing service operations on the first tool-supporting frame at the service station. Accordingly, service operations on the first tool-supporting frame and assembly operations on the workpiece with the second tool-supporting frame can be simultaneously performed."}
{"text": "Decorated food products, such as cakes, are popular items, particularly for special occasions, such as birthdays, holidays, weddings, anniversaries and other celebratory events. These items are typically purchased from bakeries since the items generally require some skill and/or apparatus to create. Currently, these items are created from colored frostings applied manually by trained bakery chefs. In addition, pre-made decorative items may be placed on a food product. These pre-made items are generally produced in bulk in a factory and inventoried until use. The use of such pre-made decorative items precludes uniquely decorated and/or personalized decorations. Also, these pre-made decorative items are generally pre-ordered, which requires a long lead time, or stored in inventory at the risk of under-ordering or over-ordering, as well as the cost of such inventory.\nAnother alternative has been to utilize an automated system for decorating cakes. One such system is disclosed in U.S. Pat. No. 5,505,775, issued to Kitos. This system utilizes an integrated work surface controlled with a computer system. An image is scanned into the computer system. The computer system then uses a motion control system to manipulate a drop on demand colorant expulsion system over a cake carried on the work surface to reproduce the scanned image. This system produces the decorations and images directly onto the food product.\nAn automated system for printing images onto edible media, that can then be applied to a food product or consumed as is, is disclosed in U.S. Pat. No. 7,286,258, issued to Schnoebeien et al. This system discloses a local processing unit and a number of image sources, including a scanner and a database of digital images stored on a local computer memory system. The user selects an image from one of the image sources and prints the image onto an edible media.\nVarious websites exist for allowing users to customize decorations. One such online website is www.genometri.com which allows users to decorate and purchase customized decorated cakes. Users can choose from pre-designed decorative templates or can choose to customize their cake decoration. Users can upload personal photographs to use in the cake decoration and can add text and clip art to the cake decoration. When the user is finished decorating the cake, the user purchases a cake decorated with the customized decoration. The purchase transaction is processed through the website."}
{"text": "In a method for detecting analyte which is well-known in this field, latex or erythrocyte is sensitized with antigen or antibody beforehand, and, then, agglutination reaction with an analyte which exists in biological fluid such as protein or a certain specific organic material is utilized for detecting the analyte. An example of latex which has widely been employed for the above-mentioned method is polystyrene latex, which is being studied also as to what influence is exerted on analyte-detecting rate or the sensitivity of detection (e.g., non-specific adsorption of materials other than analyte, or the like) by factors such as particle size and functional group (carboxyl group and the like) on the surface of particle (see Journal of the Japan Society for Clinical Laboratory Automation, 12, 2, 121-124, 1987).\nAgglutination of latex particles via analyte has a major influence on detecting rate. It takes, however, a considerable time for the agglutination, since the speed of said agglutination is determined by diffusion. Hence, it has been tried, for instance, to make latex particles carry magnetic particles and to raise the rate of agglutination or precipitation of latex particles by magnetic force (see, e.g., J. Applied Polymer Science, Vol. 50, 765-776, 1993).\nPolystyrene latex particles and the like which are mentioned in the former of the above-mentioned documents are liable to agglutinate non-specifically, and are sometimes slightly inferior in dispersion stability. Furthermore, the surface of such particles is apt to non-specifically adsorb impure protein other than analyte. Thus, a part of the inventors of this invention has provided, as a quantum dot to detect analyte in biological sample, composite particles which are composed, for instance, of semiconductor-ultra fine particles-encapsulating fine particles whose surface is covered with chain of polyethylene glycol (hereinafter referred to also as PEG) or poly(ethylene oxide) (or wherein PEG having high mobility in an aqueous solution takes plural brush-like structures) (see pamphlet of WO 02/056020, pages 16-17). Although this type of quantum dot prevents non-specific adsorption of impure protein, it has generally a particle size up to hundreds of nanometers, and has high-mobility PEG on its surface. Thus, quantum dot makes a very stable dispersion in an aqueous solution, and has therefore not necessarily been considered to be suitable for use in a system to rapidly detect analyte with use of agglutination reaction."}
{"text": "1. Field of the Invention\nThis present invention relates, generally, to the field of gaming, and more particularly to playing cards used in games.\n2. Description of Related Art\nVarious games require the use of one or more playing cards. While there are many possible variations of playing cards, a standard deck of playing cards generally consists of fifty-two playing cards, each with different markings. The fifty-two cards may be categorized by their markings into four different suits with thirteen ranks in each suit. The four suits are generally hearts, diamonds, spades, and clubs. The thirteen ranks are generally Ace (A), 2, 3, 4, 5, 6, 7, 8, 9, 10, Jack (I), Queen (Q), and King (K). The suits and ranks may have different values and/or significance in different games. Occasionally, a standard deck is packaged with one or more Joker cards.\nEach playing card in a standard deck is generally constructed of a sheet of paper or plastic. One side, the face side, displays the suit and rank. The other side, the back side, generally does not indicate the markings on the face side. This back side may instead display a variety of designs. Possible designs may include a manufacturer's logo, a distributor's design, an artistic design, a novelty design, or any of various custom designs. In a standard deck of cards, all fifty-two cards generally display the same design on their back side.\nExamples of games commonly played with at least one standard deck of playing cards includes poker, blackjack, bridge, gin rummy, go fish, and various others. Many of these card games may include elements of chance. To ensure a fair game in a game of chance, players must be dealt playing cards at random. Such randomization may be achieved through shuffling of the cards. Shuffling the playing cards randomly reorders a deck to produce a different permutation in the order of the playing cards of the deck, which contributes to the random distribution of playing cards to each player\nThere are a variety possible shuffling techniques and a wide range in the quality of randomization resulting from these techniques. Good randomization of a deck of cards may require multiple shuffles using different shuffling techniques. Some games require that a deck be randomized multiple times throughout the course of the game or series of games. These multiple acts of randomization may consume a good deal of playing time. Furthermore, shuffling skills may vary from player-to-player due to differences in experience and/or manual dexterity. There are also many opportunities for mistakes, such as inadvertently revealing the markings on one or more cards during the act of shuffling. In the event of such a mistake, the exposed card may simply be returned to the deck, which may be shuffled again or subjected to more extensive shuffling to ensure the randomization of the deck. Some games and/or the players of the same may regard such mistakes very seriously. As such, some players may choose to abort the current game and start a new one as a result of the card exposure. In, the hands of an inexperienced and/or unskilled player, shuffling may be even more time-consuming than it already needs to be and possibly even disruptive to the game.\nInexperienced and/or unskilled shuffles may also result in poorly randomized cards. Poor randomization may decrease the chance element in games of chance, which may give some players an unfair advantage or disadvantage. Poor randomization may also affect a player's ability to bluff if other players can predict what cards he or she actually holds. The ability to bluff may be a very important part of the strategy in various card games (e.g., poker). Bluffing depends, in part, on restricting players from seeing or knowing the suit or rank of the playing cards dealt to their opponents. Dealing playing cards face-down may be one way to allow a player to access and see only his or her own playing cards. Different types of playing cards, decks of playing cards, and/or card games may utilize different systems of markings, shuffling, and/or dealing.\nAnother possible danger with respect to randomization is intentional and illicit manipulation of the cards by a player who is highly skilled at shuffling and/or dealing. A player with a stake in the outcome of the game may be capable of manipulating that outcome by deliberately shuffling and/or dealing the cards into a non-randomized and/or predictable order to achieve a predetermined, unfair result. To minimize manipulation, players may perform a procedure commonly referred to as ‘cutting the deck.’ After the deck is shuffled, a non-shuffling party may ‘cut the deck’ by selecting a random number of contiguous cards from anywhere in the deck and moving those cards elsewhere in the deck, while keeping all cards face-down. This procedure, when combined with shuffling and dealing, further adds to the time-consuming nature of ensuring proper randomization of the deck of cards, especially if shuffling, dealing, and cutting have to be performed multiple times. There is, therefore, a need for improved systems and methods for efficiently maintaining the randomness and secrecy of playing card distribution."}
{"text": "The invention relates to the bonding of an electronic component to its molded package to render the package fluid-tight.\nWhile various types of protective coatings for electronic components have been known and used for various reasons for a long time, it is only recently because of changes in technology and materials that it has become critical to seal completely these units against fluid penetration.\nIn the past, protective resilient coatings have found utility in providing stress and shock relief for components mounted in hollow shells. Resilient coatings have been used also with potted units to provide protection from vibration, thermal shock, and contamination (including moisture) and from the potting or encapsulating compositions.\nProtective precoats have been used with electronic components because of differences in thermal expansion between the units and molded cases and to prevent contamination of the units by encapsulating materials.\nAs fully molded units became more popular, the seal between the leads or pins became more of a problem as the molding materials contain mold release agents to facilitate removal of the units from the molds. Unfortunately, these agents also prevent good adhesion to the metal leads or pins.\nAs water-soluble fluxes are replacing rosin fluxes in the reflow soldering of these units to circuit boards, because of environmental concerns, the sealing of the leads or pins to the covers or cases has become even more critical as these water-soluble materials are much higher in ionic material than prior rosin-based fluxes and eventually cause internal contamination and corrosion leading to failure of the unit. Thus, a need has arisen to find sealant materials which are compatible with the materials already in use in molded components and yet resist penetration by fluxes and ionic contaminants. One such solution to the above problem, described in copending application Ser. No. 863,545 filed Dec. 22, 1977, is to fill in the space between the unit (including leads) and the case with a sealant layer."}
{"text": "X-ray tubes are used in a variety of industrial and medical applications. For example, X-ray tubes are employed in medical diagnostic examination, therapeutic radiology, semiconductor fabrication, and material analysis. Regardless of the application, most X-ray tubes operate in a similar fashion. X-rays, which are high frequency electromagnetic radiation, are produced in X-ray tubes by applying an electrical current to a cathode to cause electrons to be emitted from the cathode by thermionic emission. The electrons accelerate towards and then impinge upon an anode. The distance between the cathode and the anode is generally known as A-C spacing or throw distance. When the electrons impinge upon the anode, the electrons can collide with the anode to produce X-rays. The area on the anode in which the electrons collide is generally known as a focal spot.\nX-rays can be produced through at least two mechanisms that can occur during the collision of the electrons with the anode. A first X-ray producing mechanism is referred to as X-ray fluorescence or characteristic X-ray generation. X-ray fluorescence occurs when an electron colliding with material of the anode has sufficient energy to knock an orbital electron of the anode out of an inner electron shell. Other electrons of the anode in outer electron shells fill the vacancy left in the inner electron shell. As a result of the electron of the anode moving from the outer electron shell to the inner electron shell, X-rays of a particular frequency are produced. A second X-ray producing mechanism is referred to as Bremsstrahlung. In Bremsstrahlung, electrons emitted from the cathode decelerate when deflected by nuclei of the anode. The decelerating electrons lose kinetic energy and thereby produce X-rays. The X-rays produced in Bremsstrahlung have a spectrum of frequencies. The X-rays produced through either Bremsstrahlung or X-ray fluorescence may then exit the X-ray tube to be utilized in one or more of the above-mentioned applications.\nIn certain applications, it may be beneficial to lengthen the throw length of an X-ray tube. The throw length is the distance from cathode electron emitter to the anode surface. For example, a long throw length may result in decreased back ion bombardment and evaporation of anode materials back onto the cathode. While X-ray tubes with long throw lengths may be beneficial in certain applications, a long throw length can also present difficulties. For example, as a throw length is lengthened, the electrons that accelerate towards an anode through the throw length tend to become less laminar resulting in an unacceptable focal spot on the anode. Also affected is the ability to properly focus and/or position the electron beam towards the anode target, again resulting in a less than desirable focal spot—either in terms of size, shape and/or position. When a focal spot size or location is unacceptable, it may be difficult to produce useful X-ray images.\nThe subject matter claimed herein is not limited to embodiments that solve any disadvantages or that operate only in environments such as those described above. Rather, this background is only provided to illustrate one exemplary technology area where some embodiments described herein may be practiced."}
{"text": "Ultrasonic echoscopy provides information about an examined object which may be displayed in the form of an ultrasonic echogram. Such an echogram consists of a display of acoustic impedance discontinuities or reflecting surfaces in the object. It is obtained by directing a short pulse of ultrasonic energy, typically in the 1-30 MHz frequency range, along a line called the beam axis into the examined object where any acoustic impedance discontinuities in the object reflect and return some of the energy along the same beam axis in the form of an echo. This echo is received, converted into an electric signal and displayed as an echogram on a cathode ray oscilloscope, a film, a chart or the like.\nThe echogram may constitute either a one dimensional or a two dimensional representation and in both cases the information is contained in the position and magnitude of the echo displayed. In a one dimensional display, the position along a base line is used to indicate the distance to the reflecting surface whilst the magnitude of the echo is displayed, for example, as a deflection of the base line or as an intensity change. In a two dimensional display, the position along a base line is used to indicate the distance to the reflecting surface as in a one dimensional display, and the direction of the base line is used to represent the direction of propagation of the acoustic energy which is the beam axis. The two dimensional display is obtained by changing this direction of propagation of the acoustic energy and by instituting a similar but not necessarily identical movement of the base line of the display. The magnitude of the echo is displayed as for a one dimensional display; for example, as a deflection of the base line or as an intensity change.\nThe technique of ultrasonic echoscopy is used in medical diagnosis to obtain information about the anatomy of patients. The application of this technique is now widely investigated and is described, for example, by D. E. Robinson in \"Proceedings of the Institution of Radio and Electronics Engineers Australia,\" Vol. 31, No. 11, pages 385-392, November, 1970, in his paper entitled: \"The Application of Ultrasound in Medical Diagnosis\". As pointed out in this article, ultrasonic echoscopy may be used to produce displays resembling anatomical cross-sections which have proved clinically useful when the desired information concerns physical dimensions, shapes of organs or structures or the like. Ultrasonic echography has proved of particular value as a diagnostic aid in the abdomen and pregnant uterus, eye, breast, brain, lung, kidney, liver and heart, these being areas of soft tissue with little bone and air. In general, the technique is considered to complement other techniques to provide a more complete picture of the patients condition, however particularly in pregnancies, ultrasonic echoscopy may be useful in place of X-rays where the latter may not give sufficient information or may be dangerous. In medical use, a pulse of ultrasonic energy is transmitted into a patient in a known direction and echoes are received from reflecting surfaces within the body. The time delay between a transmitted pulse and the received echo depends on the distance from the transmitter to the reflecting surface and the distance information so obtained may be displayed in a suitable way for interpretation and clinical use as a one dimensional range reading or as a two dimensional cross section as previously described.\nThis known system has sufferred from a disadvantage due to the time required to obtain a cross-sectional picture. The cross-sectional picture is made up of a multiplicity of lines of information corresponding to each beam axis position at which a pulse was transmitted and echoes received. The time required to obtain each line of information is fixed by the depth of the tissues of interest and the velocity of propagation of sound in the tissues to be examined. For a particular area of interest neither of these parameters is under the control of the operator and they form a basic limitation on the time required to obtain an echogram.\nIn U.S. Pat. No. 3,789,833 to Bom, there is disclosed the formation of an array of transducer elements arranged in a line, each providing a separate ultrasonic line of sight. Each array element is pulsed in turn, returned echoes are received and displayed on a screen and then another array element is pulsed. One primed limitation of the system disclosed by Bom is that the line spacing fixes the size of the transducer elements and the resulting ultrasonic resolution. U.S. Pat. No. 3,881,466 to Wilcox discloses the use of a plurality of transducer elements to form each ultrasonic beam, thus breaking the link between line spacing and resolution. In addition, U.S. Pat. No. 3,166,731 to Joy and 3,086,195 to Halliday disclose the application of time delays to signals associated with the various elements of a transducer element array in forming the ultrasonic beam, to cause an electronic steering and focusing action within the plane of scan.\nIt will thus be apparent that ultrasonic transducers consisting of an array of elements are being used to provide rapid cross sectional imaging, particularly in medical diagnosis. These transducers are usually rectangular piezo-electric ceramic elements sandwiched between two electrodes. The thickness of the element is selected so that the element resonates at the required megahertz ultrasonic frequency (that is, the frequency of the ultrasonic energy to be propagated into the object or medium to be examined). At megahertz frequencies, this means that the thickness of the element is approximately equal to half the wavelength. Imaging is obtained by energising either one or a group of elements to provide a single line of sight of ultrasonic information and a cross sectional image may be built up by sequentially energising the elements as described above. Because of considerations imposed by the number of lines used to form the image and lateral resolution requirements, the width of the elements forming the array typically ranges from the wavelength to two wavelengths, whilst the length dimension of the rectangular elements is usually greater than ten wavelengths.\nA common method of constructing such transducer arrays is to use a relatively large rectangular ceramic transducer, corresponding in size to the transducer array, and to divide this transducer to form the individual transducer elements of the array.\nOne method of construction includes the scribing of fine lines on one or both of the surfaces of the relatively large transducer, to cut through one or both of the electrodes to provide electric insulation between individual elements of the array. This method of array construction results in an array with high sensitivity, for the unenergised transducer material lying alongside an activated array element damps and restrains the transfer of energy from the thickness mode (in which the element is energised) into other modes which could be forced into resonance due to the mechanical coupling that exists within the massive piezo-electric ceramic forming the elements of the array. Unfortunately, the presence of this material allows a small transfer of energy between adjacent elements, which results in a significant cross talk between elements.\nIn an alternative construction, the relatively massive piezo-electric ceramic transducer is cut through completely to forming form a series of rectangular elements. This method of construction reduces the degree of cross talk between elements since, due to the complete cut, there is no direct mechanical coupling between the adjacent transducer elements. However, as the thickness and the width of the transducer elements are now comparable, there is significant transfer or coupling of energy between adjacent elements due to resonance effects. Consequently, the sensitivity of the element, and thus of the array, is reduced."}
{"text": "1. Field of Invention\nThe invention relates generally to detecting and processing acoustic signal data and more specifically to reducing noise in acoustic systems.\n2. Art Background\nAcoustic systems employ acoustic sensors such as microphones to receive audio signals. Often, these systems are used in real world environments which present desired audio and undesired audio (also referred to as noise) to a receiving microphone simultaneously. Such receiving microphones are part of a variety of systems such as a mobile phone, a handheld microphone, a hearing aid, etc. These systems often perform speech recognition processing on the received acoustic signals. Simultaneous reception of desired audio and undesired audio have a negative impact on the quality of the desired audio. Degradation of the quality of the desired audio can result in desired audio which is output to a user and is hard for the user to understand. Degraded desired audio used by an algorithm such as in speech recognition (SR) or Automatic Speech Recognition (ASR) can result in an increased error rate which can render the reconstructed speech hard to understand. Either of which presents a problem.\nUndesired audio (noise) can originate from a variety of sources, which are not the source of the desired audio. Thus, the sources of undesired audio are statistically uncorrelated with the desired audio. The sources can be of a non-stationary origin or from a stationary origin. Stationary applies to time and space where amplitude, frequency, and direction of an acoustic signal do not vary appreciably. For, example, in an automobile environment engine noise at constant speed is stationary as is road noise or wind noise, etc. In the case of a non-stationary signal, noise amplitude, frequency distribution, and direction of the acoustic signal vary as a function of time and or space. Non-stationary noise originates for example, from a car stereo, noise from a transient such as a bump, door opening or closing, conversation in the background such as chit chat in a back seat of a vehicle, etc. Stationary and non-stationary sources of undesired audio exist in office environments, concert halls, football stadiums, airplane cabins, everywhere that a user will go with an acoustic system (e.g., mobile phone, tablet computer etc. equipped with a microphone, a headset, an ear bud microphone, etc.) At times the environment the acoustic system is used in is reverberant, thereby causing the noise to reverberate within the environment, with multiple paths of undesired audio arriving at the microphone location. Either source of noise, i.e., non-stationary or stationary undesired audio, increases the error rate of speech recognition algorithms such as SR or ASR or can simply make it difficult for a system to output desired audio to a user which can be understood. All of this can present a problem.\nVarious noise cancellation approaches have been employed to reduce noise from stationary and non-stationary sources. Existing noise cancellation approaches work better in environments where the magnitude of the noise is less than the magnitude of the desired audio, e.g., in relatively low noise environments. Spectral subtraction is used to reduce noise in speech recognition algorithms and in various acoustic systems such as in hearing aids. Systems employing Spectral Subtraction do not produce acceptable error rates when used in Automatic Speech Recognition (ASR) applications when a magnitude of the undesired audio becomes large. This can present a problem.\nIn addition, existing algorithms, such as Spectral Subtraction, etc., employ non-linear treatment of an acoustic signal. Non-linear treatment of an acoustic signal results in an output that is not proportionally related to the input. Speech Recognition (SR) algorithms are developed using voice signals recorded in a quiet environment without noise. Thus, speech recognition algorithms (developed in a quiet environment without noise) produce a high error rate when non-linear distortion is introduced in the speech process through non-linear signal processing. Non-linear treatment of acoustic signals can result in non-linear distortion of the desired audio which disrupts feature extraction which is necessary for speech recognition, this results in a high error rate. All of which can present a problem.\nVarious methods have been used to try to suppress or remove undesired audio from acoustic systems, such as in Speech Recognition (SR) or Automatic Speech Recognition (ASR) applications for example. One approach is known as a Voice Activity Detector (VAD). A VAD attempts to detect when desired speech is present and when undesired speech is present. Thereby, only accepting desired speech and treating as noise by not transmitting the undesired speech. Traditional voice activity detection only works well for a single sound source or a stationary noise (undesired audio) whose magnitude is small relative to the magnitude of the desired audio. Therefore, traditional voice activity detection renders a VAD a poor performer in a noisy environment. Additionally, using a VAD to remove undesired audio does not work well when the desired audio and the undesired audio are arriving simultaneously at a receive microphone. This can present a problem.\nAcoustic systems used in noisy environments with a single microphone present a problem in that desired audio and undesired audio are received simultaneously on a single channel. Undesired audio can make the desired audio unintelligible to either a human user or to an algorithm designed to use received speech such as a Speech Recognition (SR) or an Automatic Speech Recognition (ASR) algorithm. This can present a problem. Multiple channels have been employed to address the problem of the simultaneous reception of desired and undesired audio. Thus, on one channel, desired audio and undesired audio are received and on the other channel an acoustic signal is received which also contains undesired audio and desired audio. Over time the sensitivity of the individual channels can drift which results in the undesired audio becoming unbalanced between the channels. Drifting channel sensitivities can lead to inaccurate removal of undesired audio from desired audio. Non-linear distortion of the original desired audio signal can result from processing acoustic signals obtained from channels whose sensitivities drift over time. This can present a problem."}
{"text": "Cryptophycins are secondary metabolites belonging to the class of depsipeptide macrocycles produced by cyanobacteria of the genus Nostoc. Their name refers to the fact that they are highly cytotoxic towards yeasts of the genus Cryptococcus. The first representative of this class of molecules, cryptophycin-1 (C-1), was isolated in 1990 from cyanobacterium Nostoc sp (ATCC 53789) (see Eiβler S., et al., Synthesis 2006, 22, 3747-3789). The structure, the general formula and the numbering of the carbon atoms of these compounds, as described in WO 98/08505, are recalled below:\n\nCryptophycins C1 and C-52, which are characterized by an epoxide function represented below, have anticancer properties. Phase II clinical trials in lung cancer were conducted with C-52 (LY 355073): see Edelman M. J., et al., Lung Cancer 2003, 39, 197-199; Sessa C., et al., Eur. J. Cancer 2002, 38, 2388-96. Cryptophycin C-55, a prodrug of C-52, is itself characterized by a chlorohydrin function instead of the epoxide function (Bionpally R. R., et al., Cancer Chemother Pharmacol 2003, 52, 25-33). C-55 proved to be very active, but is not stable in solution. Derivatives of chlorohydrin glycinate type such as the compound C-55 Gly have also been described as gaining in stability (Liang J., et al., Investigational New Drugs 2005, 23, 213-224).\n\nConjugate chemistry has been known for many years and has been applied to several families of cytotoxic agents, for instance the maytansinoids (WO 04/103 272), taxanes (WO 06/061 258), tomaymycins (WO 09/016,516), leptomycins (WO 07/144,709), CC-1065 and analogues thereof (WO 2007/102 069); see also, with regard to conjugates, Monneret C., et al., Bulletin du Cancer. 2000, 87(11), 829-38; Ricart A. D., et al., Nature Clinical Practice Oncology 2007, 4, 245-255. However, it has not been applied to cryptophycin derivatives conjugated to antibodies or to other targeting agents.\nThe technical problem that the present invention intends to solve is that of proposing novel conjugates based on cryptophycin derivatives, and also novel cryptophycin derivatives that are suitable for being conjugated.\nEP 0 830 136 and WO 98/08505 describe cryptophycin derivatives but do not describe cryptophycin conjugates. WO 98/08505 describes cryptophycin derivatives of formula (A):\nin which Ar may represent a group Ar′ of formula:\nin which R54 represents H, a group (C1-C6)alkyl, (C1-C6)alkyl(R57, R57′, R57″), aryl, phenyl, heterocycloalkyl, a halogen atom, COOR57, PO3H, SO3H, SO2R58, NR59R60, NHOR61, NHOR61′, CN, NO2, OR62, CH2OR62′, CH2NR96R96′, (C1-C6)alkylOR100,\nSR63; R55 and R56 represent H, a group (C1-C6)alkyl, C(R57R57′R57″), aryl, phenyl, heterocycloalkyl, a halogen atom, COOR57, PO3H, SO3H, SO2R58, NR59R60, NHOR61, NHCHR61′, CN, NO2, OR62, CH2OR62′, CH2OCOR95, CH2NR96R96′, (C1-C6)alkylOR100, (C1-C6)alkylNR59R60. WO 98/08505 especially describes the following compounds:\nwhich are characterized by the end group —NHC(═O)Ot-Bu. WO 98/08505 does not specify that these compounds are suitable for or intended to be conjugated.\nUS 2007/0 213 511 describes calicheamicin immunoconjugates. Among these, mention is made of Mylotarg® (or CMA-676), which is a calicheamicin immunoconjugate used in the treatment of AML (anti-CD33-calicheamicin). See also in respect of conjugates: WO 2006/042 240.\nAl-awar R. S., et al., J. Med. Chem. 2003, 46(14), 2985-3007 and Al-awar R. S., et al., Mol. Cancer Ther. 2004, 3(9), 1061-1067 describe cryptophycin derivatives, and also their in vivo evaluations.\nWO 08/010,101 describes an anti-EphA2 monoclonal antibody and also the corresponding conjugates comprising one or more molecules of a cytotoxic compound attached to the monoclonal antibody. WO 08/047,242 describes an anti-CD38 monoclonal antibody and also the corresponding conjugates comprising one or more molecules of a cytotoxic compound attached to the monoclonal antibody. The cytotoxic compound may be chosen from maytansinoids, taxanes, tomaymycins, leptomycins and CC-1065 and analogues thereof.\nWO 2009/126 934 describes anti-CD70 antibodies and their conjugates with cytotoxic compounds; cryptophycin is mentioned among the cytotoxic agents. WO 2009/134 976 describes conjugates with an optimized degree of substitution for delivering the required amount of cytotoxic agent into the cell; cryptophycin is mentioned among the cytotoxic agents.\nWO 2005/116 255 describes conjugates of aptamers and of a cytotoxic agent that may be a cryptophycin (see [0037] and Table 2), the linker possibly comprising a PEG chain ([0038]). More particularly, the cryptophycin Cryp-NH2 is described:\nand also its conjugates of aptamers with the following linkers: NHS-PEG-erythritol, pNP-PEG-erythritol, NHS-PEG-octaPEG, pNP-PEG-octaPEG and PEG-comb (Tables 3 and 4). The nature of the sequence on the phenyl ring (—CH2O—C(═O)—CMe2-CH2NH— . . . ) is different from that which is envisaged in the present invention. WO 2006/096 754 also describes conjugates of aptamers and of a cytotoxic agent that may also be a cryptophycin.\nWO 2009/002 993 describes cytotoxic conjugates of formula B-L-A comprising hydrophilic linkers, for example the linker of formula:\n\nThe cytotoxic agent may be a cryptophycin (page 46), but the point of attachment of the linker is not specified. An example of a conjugate is EC0262:\nwhose linker and point of attachment are different from those envisaged in the present invention."}
{"text": "A memory module typically includes a small printed circuit board which is provided with memory chips, such as dynamic random access memories (DRAMs). Memory modules typically form or extend, for example, the working memory of electronic devices, such as computers or printers. A special type of memory module is a dual inline memory module (DIMM) which includes two rows of contact pins, one on the front side and one on the rear side of the printed circuit board respectively.\nOne or several memory modules are typically controlled by a memory controller. In order to design this memory controller, knowledge of the longest distance of a line (i.e., a trace) between the memory controller and a memory chip which is to be controlled by the memory controller, which is generally located inside a memory module, is of interest by reason of the running time which is to be taken into consideration by the memory controller. Accordingly, a distinction is made between: a short line (e.g., approximately 5 cm), such as for graphics applications; a medium length line (e.g., approximately 20 cm), such as for personal computers (PCs); and, a long line of more than, for example, 30 cm, such as for connections in a backplane.\nSince the reuse of components plays a significant role in today's semiconductor industry, most memory controllers are currently designed for a medium line length, although assemblies which by reason of their configuration could operate with a memory controller designed for a short line length, would be advantageous by reason of their shorter running time and thus higher clock frequency of the memory controller with respect to a memory controller for a medium length line.\nHowever, the length of the line is also typically determined in particular by the dimensions and designs of the memory modules."}
{"text": "1. Technical Field\nThe present invention relates generally to reacting a fluid with a body, and more particularly to a body for reacting with a fluid.\n2. Description of Related Art\nFluids (e.g., gases, liquids and the like) may be passed over or through a variety of substances in order to facilitate a reaction. The reaction of species in the fluid, the removal of species from the fluid, and the conversion of the fluid in some chemical or physical fashion may be enhanced by contacting the fluid with a suitable body. In some cases, reactions occur at the interface between the fluid and a surface of the body, and in such cases, it may be advantageous to maximize this surface area (e.g., by using a porous body). Reactions between a fluid and a body may also be used to modify the body itself.\nMany suitable bodies are substantially porous (e.g., over 5%, 10%, 20%, 40%, 60%, 80%, or even over 90% porous). Porosity may be continuous, and in such cases, the fluid may substantially saturate the body. When pushed by a suitable driving force (e.g., pressure, voltage, magnetic field or other gradient in thermodynamic potential) a fluid may be caused to pass through the body. In some cases, a species may be removed from a fluid by passing the fluid through pores that block the passage of the species (e.g., filtration). In other cases, a species may be removed or reacted to form another species (e.g., in heterogeneous catalysis) and/or combine with the body per se (e.g., gettering). A species may also be substantially adsorbed or absorbed by the body.\nMany bodies have been fabricated by choosing relatively pure starting materials and mixing them, and increasingly complex bodies require the addition of more starting materials. For example, a catalytic converter substrate may be required to have certain thermal, mechanical, and surface area properties, so might be made of cordierite (Mg2Al4Si5O18), which may be made by mixing MgO, Al2O3 and SiO2 starting materials. Cordierite containing iron (Mg,Fe)2Al4Si5O18 may have some improved properties, and might be made by adding FeO or Fe2O3 to the aforementioned mixture. Sintering aids, grain boundary phases, and catalytic species could be similarly added, and improvements to many materials generally entail the addition of further components.\nMany useful bodies, particularly bodies for high temperature applications, include several components, and in some cases, additional components may improve properties. For example, mullite (3Al2O3-2SiO2) materials often have high strength at fairly high temperatures (e.g., 1300 C), and U.S. Pat. No. 3,954,672 identifies certain cordierite-mullite compositions (i.e., of increased complexity) that have some improved properties over mullite. As such, materials having generally improved properties in an application may often be more complex than known materials typically used in the application.\nMany useful bodies are fabricated from combinations of SiO2, Al2O3, FeO, MgO, CaO, TiO2 and other materials, and often include one or more useful phases (e.g, mullite, cordierite, spinels, perovskites, amorphous), each of which may include several components. Thus, the discovery of improved bodies for a variety of applications might be enhanced by basing those bodies on compositions known to have useful properties, then increasing complexity around these compositions.\nPorous bodies may be used for filtration, including without limitation deep bed filtration, wall flow filtration, cake filtration, and other methods. Generally, an appropriate body for use in filtration may be chosen based upon a variety of factors, including required flow rates, viscosity of the fluid, phase assemblage of the fluid (e.g., suspended solids in a liquid, emulsions), concentration of species (to be treated) in the fluid, desired pressure differential (if pressure is driving the fluid through/past a body), temperature of the application, chemical reactivity (or lack thereof) and other factors. Available geometrical and mass constraints may also determine an appropriate filtration method. For example, large “ponds” of deep bed filtration bodies may be used to filter large amounts of wastewater, whereas catalytic removal of contaminants in an automotive exhaust gas stream may require a small, portable body.\nIn some applications, the mechanical behavior of the body may be important. Often, the driving force used to cause a fluid to pass through or past a body creates a mechanical stress in the body itself, and the body's resistance to this mechanical stress is a requirement in many applications. Some applications require that a body have sufficient mechanical strength to withstand an applied pressure exerted by the fluid (e.g., in a filter). Some applications may require a low thermal expansion coefficient (CTE), good thermal shock resistance, or good thermal shock damage resistance.\nIn many applications, channels or other substantially “open” regions of the body (offering minimal impedance to fluid flow) are used to substantially increase area for reaction or filtration (for example, as in U.S. Pat. Nos. 4,253,992 and 4,276,071). In such applications, relatively thin walls separate regions having substantially minimal impedance to fluid flow. Walls separating the channels should have both high porosity to maximize surface area or permeability, but not so high porosity that mechanical properties are degraded, and the pore size distribution should provide for the desired treatment (e.g., actually filter the species being removed).\nIn some applications (e.g., a filter bed) a body may be essentially hydrostatically supported during operation, and so require little mechanical strength (e.g., shear or tensile strength) during filtration. Some applications also include backwashing, which often creates a different mechanical stress than that created during filtration. In such instances, some mechanical strength or appropriate containment may be necessary. Thermal stresses, thermal expansion mismatch, changes in crystallographic structure, physical impact, and other factors may also create certain requirements of a body in a given application.\nIncreasingly, cost may be an important factor in a given application. Costs may include capital costs associated with fabricating the body and associated fluid control system itself. Cost may also include operational costs. Costs may also include disposal costs, environmental costs, “cradle to grave” costs and other costs associated with implementing a particular treatment solution. The energy required to create and implement a particular body may be an important cost factor, and in such cases, reducing the energy required to make and use a particular device may be advantageous. Cost may include a cost associated with emitting global warming gases, environmental pollutants, or other contaminants. Often, minimizing the embodied energy associated with a product (e.g., the energy required to create and implement the product) and/or minimizing a total lifecycle cost of the product may be advantageous. The implementation associated with a treatment method, the method of treating the fluid with the body, the disposal of treated substances and/or the body itself, and other lifecycle costs may generally include both capital costs (including raw materials costs), operational costs (including lifetime) and disposal/recycling costs."}
{"text": "This invention relates to flameless atomic absorption spectroscopy, and more particularly the invention relates to a graphite tube assembly for use in measuring the flameless atomic absorption, particularly of liquid samples.\nFor atomic absorption measurements, it is conventional for a sample to be introduced into the interior of the graphite tube, dried at a relatively low temperature, ashed and, subsequently, atomized at a high temperature. The graphite tube is placed in the path of the rays of a spectrometer so that the absorption of light by the atoms present in the interior of the graphite tube can be measured.\nIt has been found that in measuring the atomic absorption of liquid samples disturbances may occur because the liquids spread out over wide portions of the inner wall of the graphite tube. The liquids also flow to the end portions of the graphite tube where there is incomplete evaporation, because of the relatively lower temperature thereat, so that the sample material is carried over and interferes with subsequent measurements of further samples. Because of the slightly porous structure of graphite, the sample may infiltrate the wall of the graphite tube to such an extent that an uncontrollable portion thereof will pass through the wall to the exterior when the tube is heated, and is thereby lost for the measurement. This results in substantially lowering the sensitivity of the measurement.\nIt is known, by virtue of German Offenlegungsschrift 23 23 774, that to avoid such problems the graphite tube may be provided with grooves, at least over a portion of its inner wall and extending transversely with respect to the tube axis. However, since the cutting of the grooves causes a substantial roughening of the inner face, additional infiltration of the sample material results therefrom, so that the improvement achieved does not reach the desired degree.\nFurther, it is known to manufacture graphite tubes of the aforementioned type from so-called pyrolytic graphite instead of from the usual graphite. This pyrolytic graphite is not porous and, therefore, is impervious to gases. Grooves may also be cut in the pyrolytic graphite without the sample infiltrating the walls of the tube. However, the use of pyrolytic graphite for this purpose is undesirable because of the high cost thereof."}
{"text": "Pateamine A (patA) was first isolated from the marine sponge Mycale found off the shores of New Zealand. The natural form bears a thiazole and an E,Z-dienoate within a 19-membered macrocycle and a trienylamine side chain. Two additional pateamines, pateamines B and C, were also isolated. Their structures differ from pateamine A only in the nature of the terminal group of the trienylamine side chain. The basic structure for these three isolated natural forms is shown below.\n\nIsolated pateamine A (native pateamine A) is a novel marine product that promises to be quite useful as a biochemical probe and which displays potent immunosuppressive properties with low cytotoxicity. In MLR (mixed lymphocyte reaction) assay, IC50=2.6 nM while the LCV (lymphocyte viability assay)/MLR ratio is >1000. In comparing native pateamine A to cyclosporin A in a mouse skin graft rejection assay, native pateamine A resulted in a 15-day survival period as opposed to cyclosporin A having only a 10-day survival of the skin graft. Additionally, at high doses, toxicity was at 17% in these studies. For other dose levels, there was no toxicity. All doses were active.\nIt has also been found that native pateamine A specifically inhibits an intracellular step of the T-cell receptor signal transduction pathway leading to IL-2 transcription. Two syntheses of native pateamine A have been reported. The utility of these molecules as an immunosuppressant or immunostimulant is severely restricted because the molecule lacks stability. Additionally, natural sources of the molecule are limited. Thus, continuous development of synthetic pateamine derivatives having the same or lower toxicity, potent activity and increased stability is required.\nPotent activity of native pateamine A in a mixed lymphocyte reaction and also in a mouse skin graft rejection assay have been observed. Native pateamine A originally showed activity in a mixed lymphocyte reaction, (IC50 2.6 nM) and in the mouse skin graft rejection assay. Native pateamine A was found to be more potent than cyclosporin A with only low toxicity at high doses but all doses were active. Native pateamine A may also inhibit a specific intracellular signaling pathway involved in T cell receptor-mediated IL-2 production. In addition to its effect on TCR signaling pathway, native pateamine A has been found to induce apoptosis in certain mammalian cell lines, especially those that are transformed with the oncogene Ras.\nAnalysis of native pateamine A structure reveals a rigid eastern half (C6-C24) including the thiazole, dienoate, and the triene sidechain, due to extended conjugation, and a more flexible western half (C1-C5). Furthermore, C3-Boc-PatA was found to have only 3-4 fold lower activity than native pateamine A.\nNatural products have proven to be extremely useful as probes of biological processes. Examples include the immunosuppressive, microbial secondary metabolites, cyclosporin A, FK506, and rapamycin. Marine organisms have also been a rich source of bioactive compounds, which are proving useful as drug leads and biological probes. For example, bryostatin, epithilone, discodermolide and ecteinascidin show great potential as anti-cancer agents and have revealed novel biological mechanisms of action.\nImmunomodulators are useful for treating systemic autoimmune diseases, such as lupus erythematosus and diabetes, as well as immunodeficiency diseases. Immunomodulators are also useful for immunotherapy of cancer or to prevent rejections of foreign organs or other tissues in transplants, e.g., kidney, heart, or bone marrow. Examples of immunomodulators include: FK506, muramylic acid dipeptide derivatives, levamisole, niridazole, oxysuran, flagyl, and others from the groups of interferons, interleukins, leukotrienes, corticosteroids, and cyclosporins. Many of these compounds, however, have undesirable side effects and/or high toxicity in a subject in need thereof. New immunomodulator compounds are, therefore, needed to provide a wider range of immunomodulator function for specific areas with a minimum of undesirable side effects.\nUnfortunately, many immunomodulators available currently have undesirable side effects and/or high toxicity and are often difficult to synthesize in pharmacologically effective amounts. What is needed is one or more immunomodulative compounds that may be synthetically produced in effective amounts and provide a broad spectrum of activity (e.g., immunomodulator function) with increased stability and with less undesirable side effects."}
{"text": "1. Field of the Invention\nThis invention relates to integrated circuit functional testing and, more particularly, to system level functional testing of microprocessors.\n2. Description of the Related Art\nSemiconductor integrated circuit (IC) devices are manufactured using many complex processing steps. In many cases, more complex ICs may require more processing steps. During any of which one or more defects may be introduced into the IC. For example, microscopic airborne particulates that find their way onto the device during manufacture may cause either latent or hard defects. Many ICs are typically manufactured on one wafer (i.e., semiconductor substrate). To reduce manufacturing defects and improve product quality and yield many manufacturers have implemented ongoing quality control in the wafer fabrication and device assembly facilities.\nHowever, there are still a fair number of IC s that are either defective from the start or may have latent defects. Thus, device manufacturers typically test the ICs using various test strategies during the various manufacturing phases. For example, to screen out bad ICs prior to assembling the ICs into device packages, ICs may be tested while still part of the wafer. This type of testing is sometimes referred to as wafer-level test or wafer probe. In addition, since defects may be introduced during the assembly process, each IC may be tested again, after assembly, during what is sometimes referred to as production or final testing.\nProduction testing may take many forms depending on the various factors including customer requirements. The goal is to screen out all devices that have defects. In some cases, the ICs may be operated under various operational and environmental stresses to allow latent defects to manifest.\nThe type of testing performed during wafer-level and production final test typically includes the use of testing machines designed to provide necessary voltages and stimuli to the inputs of an IC. During a type of testing referred to as Functional testing, the voltages and stimuli attempts to place the device in as many operational modes as possible thereby simulating device operation. While the device is operating, the IC outputs are monitored and compared against a theoretically known good output signature that is typically stored in a test file in the tester memory. Functional testing uses a large number of test files (also referred to as test patterns) that define the stimulus and response patterns. Functional testing may generally be expected to screen a large number of defects in many fault classes.\nAn alternative and sometimes complementary test method used to screen defects is referred to as scan testing. Scan testing also uses test patterns. However, unlike functional testing, scan testing doesn't attempt to operate the device in all of its operational modes. Instead scan testing relies on test logic and testability features that have been designed into the IC that enable scan chains consisting of many clocked circuit elements such as flip-flops and some latches, for example, to be loaded with special test patterns. The scan chains may loaded in one mode (e.g., scan) and clocked in another mode (e.g., run). The test patterns are propagated through the IC logic. The scan patterns are clocked out of the device and compared to response patterns stored within tester memory. For certain fault classes such as stuck at faults, for example, scan testing has been known to obtain fault coverages as high as 99% and in some cases 100%.\nHowever, regardless of the test methods used, as IC operating speeds have increased, it has been increasingly difficult to develop testing machines that can adequately operate at the required speeds. This is especially true for microprocessor testing. Testing machines that can operate at the required frequencies are very expensive. Furthermore, neither functional testing nor scan testing may screen all defects in any fault class without using enormous test patterns requiring long test durations. To obtain very high fault coverage for all fault classes, combinations of test methods may be used. For example, due to the parallel nature of superscalar microprocessor core logic, certain types of failures, such as certain speed paths, may not necessarily manifest when an IC is run through a conventional production test flow."}
{"text": "Modern computer systems typically consist of a CPU to process data, a networking interface to communicate to other computer systems, and one or more durable storage units. The system may stop processing due to power failure, program incorrectness, or a hardware fault. Such failures are often called process failures. The durable storage units are able to keep the data intact while the fault is repaired.\nA set of these computer systems can be networked to form a cluster. Although the network is generally reliable, occasional faults may occur to disrupt communication between certain nodes or sets of nodes. This disruption in communication is often called a network partition.\nEach of these nodes runs a transactional storage system that both reads and writes data (a database management system). Some of this data is concurrently accessed by applications operating on different nodes. To guarantee data consistency, transactional database replication techniques are used to manage and regulate access to that data. However, such conventional replication techniques are associated with a number of problems.\nFor instance, conventional replication systems typically require an administrator to stop all applications, then stop the databases running at each of the nodes in the replicating group, then go to each node and notify it about the change in membership (usually by editing a table), then restarting the databases on each of these machines, then restarting the applications. This necessity of having to stop the databases running in the system can be undesirable.\nWhat is needed, therefore, are transactional database replication techniques that do not require stopping of databases during replication."}
{"text": "Investigators have developed a variety of technical approaches to the generation of coordinate pairs of signals from electrographic devices. Industrial requirements for these devices are increasing concomitantly with the evolution of computer graphics, computer-aided design, and computer-aided manufacturing systems. Thus, a considerable degree of accuracy in pinpointing physical positions upon the surfaces of the digitizers is required for many applications. Other applications of the electrographic services include touch screen devices wherein the operator's finger or a stylus or the like is used to touch a portion of the accessing surface such that it emulates a key of a keyboard.\nThe operation of a digitizer or graphics tablet generally involves the same manual procedures as are employed in conventional graphics design, a stylus or tracer representing a writing instrument being drawn across or selectively positioned upon the position responsive surface of the digitizer. In turn, the electrographic device responds to the position of the stylus to generate paired analog coordinate signals which are digitized and conveyed to a host computer facility.\nEarly approaches to digitizer structures, for example, have looked to employing compostie structures wherein a grid formed of two spaced arrays of mutually orthogonally disposed fine wires is embedded in an insulative carrier. One surface of this structure serves to yieldably receive a stylus input which yielding causes the grid to read out coordinate signals. More recent and improved approaches to achieving readouts have been accomplished through resort to a capacitive coupling of the stylus or locating instrument with the position responsive surface to generate the paired analog coordinate signals. Such capacitive coupling can be carried out either with a grid layer which is formed of spaced linear arrays of conductors or through resort to the use of an electrically resistive material layer or coating.\nAn immediately apparent advantage of developing position responsive surfaces or digitizers having writing surfaces formed of a continuous resistive material resides in the inherent simplicity of merely providing a resistive surface upon a supportive substrate such as glass or plastic. Further, unlike conventionally encountered grid structures, the resistive coatings as well as their supportive substrates may be transparent to considerably broaden the industrial applications for the devices. For example, the digitizers may be placed over graphics or photographic material for the purpose of tracing various profiles.\nA variety of technical problems have been encountered in the development of resistive coating type digitizer devices, one of which concerns the non-uniform nature of the coordinate readouts achieved with the surfaces. Generally, precise one-to-one coorespondence or linearity is required between the actual stylus or tracer position and the resultant coordinate signals. Because the resistive coatings cannot be practically developed without local resistance (thickness) variations, for example of about .+-.10%, the non-linear aspects of the otherwise promising approach has required a considerable amount of investigation and development. Exemplary of such development is the border treatment or switching technique of Turner in U.S. Pat. No. 3,699,439 entitled \"Electrical Probe-Position Responsive Apparatus and Method\" issued Oct. 17, 1972, and assigned in common herewith. This approach uses a direct current form of input to the resistive surface from a hand-held stylus, the tip of which is physically applied to the resistive surface. Schlosser et al describes still another improvement wherein an a.c. input signal is utilized in conjunction with the devices and signal treatment of the resulting coordinate pair output signal is considerably improved. See U.S. Pat. No. 4,456,787 entitled \"Electrographic System and Method\", issued June 26, 1984, also assigned in common herewith. Position responsive performance of the resistive layer devices further has been improved by a voltage waveform zero crossing approach and an arrangement wherein a.c. signals are applied to the resistive layer itself to be detected by a stylus or tracer as described in U.S. Pat. No. 4,055,726 by Turner et al. entitled \"Electrical Position Resolving by Zero-Crossing Delay\" issued Oct. 25, 1977, and also assigned in common herewith. Substantially improved accuracies for the resistive surface type digitizer devices have been achieved through a correction procedure wherein memory retained correction data are employed with the digitizer such that any given pair of coordinate signals are corrected in accordance with data collected with respect to each digitizer resistive surface unit during the manufacture of the digitizers themselves. With such an arrangement, the speed of correction is made practical and the accuracy of the devices is significantly improved. The correction table improvements for these services is described, for example, in application for United States patent, Ser. No. 664,980, filed Oct. 26, 1984, by Nakamura et al. and assigned in common herewith as well as in application for United States Pat. Ser. No. 742,733, entitled \"Electrographic System and Method,\" filed June 7, 1985, by Nakamura et al. and assigned in common herewith.\nCapacitive coupling using a stylus or locating device has been employed with grid layers which are formed as adjacent but spaced-apart arrays of elongate thin conductors. For example, these grid conductors may be provided as lines of silver ink deposited in orthogonally disposed relationships upon the opposite faces of a sheet of insulative material such as Mylar. As described in Rodgers et al., U.S. Pat. No. 4,492,819, issued Jan. 8, 1985, this grid surface may be employed with a stylus which injects an a.c. signal capacitively at the surface thereof. To detect this signal, a ladder form of resistance network is employed with each of the conductor arrays such that a predetermined resistance is coupled between each conductor from first to last and a discrete resistor is coupled from the union of two successive resistors to ground. Generally these devices operate in a current mode such that current values are determined at the peripherally disposed resistor strings. As is apparent, because of the necessity of employing a conductive form of grid line, these devices are limited to opaque position responsive surface applications. While discrete matched resistors are required to couple the grid line nodes, a conductive or carbon loaded ink advantageously may be employed to provide the grid-to-grid resistive components of the unit, however, at the expense of a resistance value deviation for each discrete increment of resistance between adjacent parallel grid lines.\nInvestigations have determined that there are operational trade-offs occasioned with the various design approaches to grid-type digitizers. For example, where a.c. signals are injected from the stylus into a passive orthogonal grid, the grid structured surface electrically appears as a high impedance to the stylus or pick-up. Thus, the presence of moisture at the tablet or digitizer surface will cause severe inaccuracies. Further, this type device is prone to react adversely in terms of readout accuracy to slightly conductive materials including certain forms of paper.\nOn the other hand where the grid-type digitizers are excited by a.c. signals emanating from the peripheral resistor strings or multi-nodal stripes, a potential gradient is established for coordinate identification. The stylus or pick-up, in turn, receives analog coordinate signals in an arrangement desirably immune from hand effects, moisture effects and the like. However, the linearity of these devices becomes severely impaired to essentially negate the advantages otherwise sought.\nWhere it is desired to provide the above-noted grid structures in a transparent embodiment for a digitizer, limitations may be observed. To achieve transparency of the grid lines themselves, a transparent material must be employed, for example an indium tin oxide. These materials, however, may exhibit an impedance or resistivity which establishes a finite resistance between the oppositely disposed ends of a tablet or position responsive surface. Such finite resistance, unless accommodated for, will impose debilitating error unless somehow corrected. The finite resistance further may impose ambiguities within the system such that a voltage at one position along one trace or grid line will occur in equal value at a different location in another grid line. A practical implementation, therefore, of grid type tablets requires a minimization of the chance of occurrence of such ambiguities."}
{"text": "1. Field of the Invention\nThe present invention relates to an image processing system and an image processing method in which an image forming apparatus is connected with an information processing apparatus via a network.\n2. Description of the Related Art\nIn recent years, multifunction peripherals in each of which a scanner function and a printer function are implemented inside one housing, an image is read using the scanner function, and the read image is formed on a sheet of paper using the printer function has wide spread. A system in which the multifunctional peripheral is connected with a computer via a network, and setting for reading by the scanner function or setting for printing by the printer function in the multifunction peripheral is performed using the computer is being developed.\nIn this system, a technique of determining a layout of image data (contents), which is read from an original document by the scanner function of the multifunction peripheral, on a sheet of paper through the computer connected with the multifunction peripheral via the network has been proposed.\nThat is, the image data obtained by reading an image through the scanner function in the multifunction peripheral is transmitted to the computer via the network, and the layout on a sheet of paper is determined at the computer side by using a layout technique. As an algorithm for arranging a plurality of contents having different sizes on a sheet of paper (a page), for example, a first-fit decreasing-height (FFDH) disclosed in E. G. Coffman, JR., M. R. Garey, D. S. Johonson, R. E. Tarjan (Bell Lab. & Stanford Univ.) “Performance bounds for level-oriented two-dimensional packing algorithms,” SIAM J. Comput., vol. 9, pp 808-826, No. 4, 1980, may be used. The computer generates drawing data using, for example, a page description language (PDL) according to the determined layout and transmits the drawing data to the multifunction peripheral via the network. The multifunction peripheral performs image formation based on the received drawing data.\nIn the conventional system described above, the image data obtained by reading the original document through the scanner function of the multifunction peripheral is transmitted, and the process of extracting contents from the received image data or the process of determining the layout of the extracted contents is performed at the compute side. That is, the computer deals with the image data whose size corresponds to the size of the original document and whose resolution is the read resolution that is set by the scanner function. Therefore, there is a problem in that a load of the computer side increases."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of in-flight entertainment systems. More particularly, the present invention relates to a seat electronics unit and corresponding replacement technique designed to reduce the amount of time required for removal and/or installation of the seat electronics unit.\n2. Description of Art Related to the Invention\nOver the last few decades, commercial aircraft has become a necessary mode of travel for personal and business reasons. In order to improve passenger comfort, many commercial airlines now offer on-board telephony as well as in-flight entertainment such as video games, pay-per-view movies, hotel reservations services, and the like. These telephony and entertainment services are controlled by a plurality of seat electronics units \"SEUs\" as shown in FIG. 1.\nTypically, a first SEU (\"SEU.sub.1 \") 10 is used to control the transfer of information between a telephone handset 30 and telephony transmitter/receiver circuitry (not shown) in order to establish and maintain audio-based communications with a telephone remotely located from the aircraft. A second SEU (\"SEU.sub.2 \") 20 is connected to a number of hand-held control units (e.g., control unit 40) and their corresponding liquid crystal displays (\"LCD\") 50, mounted onto the back of one of a group of passenger seats 60 as shown, or to an armrest of the passenger seat. SEU.sub.2 20 receives control signals from each handheld control unit of passenger seats 70 to signal other electronic devices to provide video content (e.g., movies) for display on its associated LCD or to select a menu representing one of a number of in-flight entertainment activities. Both of these SEUs 10 and 20 include electrical circuitry encapsulated within its protective cover.\nReferring to FIG. 2, each SEU (SEU.sub.2 20 for example) is connected to a pair of legs 75 of one of the passenger seats 70 by mounting hardware. The mounting hardware includes a mounting plate 100 with a number of pre-drilled mounting holes 110 arranged in a predetermined pattern. This pattern enables the mounting plate 100 to support a variety of SEU sizes provided by different original equipment manufacturers (\"OEMs\"). Bolts 120 are inserted through a flange element 125 of the SEU and mounting holes 110 of the mounting plate, and subsequently secured by hexagonal nuts (not shown). As a result, the SEU is secured to the mounting plate 100 and the flange element 125, adjacent to a first lateral side 121 of a protective cover 120 of the SEU, rests flush against the mounting plate 100.\nA plurality of communication connectors 130.sub.1 -130.sub.m (\"m\" being a positive whole number), normally multiple-pin female connectors, are implemented on a back side 122 of the protective cover 120. Collectively, as shown, connectors 130.sub.1 -130.sub.6 provide an electrical interface so that information can be transferred into or from the SEU via communication cables 140.sub.1 -140.sub.6. Of course, the communication cables 140.sub.1 -140.sub.6 are mechanically and electrically adaptive to connectors 130.sub.1 -130.sub.6.\nOver the last few years, it has been determined that the conventional SEU architecture possesses a number of disadvantages. One disadvantage involves design inflexibility in which the mounting plate is incapable of accommodating a variety of seat designs and SEU sizes. Another disadvantage is that it is time consuming to remove a conventional SEU for any number of reasons (e.g., scheduled servicing, repair after becoming completely or partially nonfunctional, etc.) as well as to install a SEU. More specifically, in order to remove the SEU, the communication cables 140.sub.1 -140.sub.6 would be initially disconnected from the SEU. After disconnecting the communication cables 140.sub.1 -140.sub.6, the bolts 120 are loosened and removed from the mounting holes 110 of the mounting plate 100. This allows the SEU to be disconnected from the mounting plate 100 and subsequently substituted with another SEU or repaired and reinstalled. Conversely, in order to install a new or repaired SEU, it is re-connected to the mounting plate 100 by re-attaching the bolts 120 and the communication cables 140.sub.1 -140.sub.6.\nThe replacement technique associated with the conventional SEU architecture clearly is labor intensive. This creates havoc with flight schedules because, if a few SEUs need to be replaced, it may result in aircraft downtime and subject commercial airline passengers to long flight delays. Such delays further increase maintenance costs creating an adverse tradeoff between maintenance costs and customer satisfaction. This tradeoff would pose adverse financial effects on the commercial airline. Moreover, this replacement technique may indirectly cause additional installation problems if the repair person accidentally damages a connector upon disconnecting and re-connecting its communication cable, mistakenly leaves the cables unconnected or improperly connected, and other related problems.\nTherefore, it would be advantageous to develop a SEU and corresponding replacement technique which overcomes the above-identified disadvantages."}
{"text": "In order to perform cipher communication in a secret key cipher system, both parties which make communication are required to hold in common a key in advance. Hitherto, sharing of the key has been performed by a system as shown in FIG. 11 for example. Referring to FIG. 11, numeral 41 designates a coder which is used by a sending party of a message (hereinafter is simply referred to as a sending party), and 42 designates a decoder which is used by a receiving party of the message (hereinafter is simply referred to as a receiving party), and the coder 41 comprises a random number generating means 43, a first enciphering means 44 and a second enciphering means 45, and the decoder 42 comprises a first decoding means 46 and a second decoding means 47.\nOperation of the conventional common ownership system of the key composed like this is described hereafter. Since the key for enciphering a plain text of a message (hereinafter is simply referred to as a plain text) m is required to be changed frequently from the aspect of safety, a random number which is generated by the random number generating means 43 is used. Hereinafter, this is described as a session key. The sending party sends an output r1 (it is called a session key as another name) of the random number generating means 43 which is held in the coder 41 to the receiving party in order to hold in common with the receiving party; but if r1 is sent as the sate of raw data it is liable to be tapped on the communication line between the coder 41 and the decoder 42, and hence 41 is enciphered by the enciphering means 44 and is sent. The key which is used to encipher r1 is called as a master key km, and it is the key which is held in common by the sensing party and receiving party in advance. The master key is used only when the session key is sent by enciphering, and is generally fixed during a long time period. The decoder 42 of the receiving party decodes the enciphered random number by the master key km and restores r1 by using the first decoding means 46. Thereby, since the sending party and receiving party have held in common the session key r1, thereafter, cipher communication of the plain text m can be accomplished by the session key r1 by using the second enciphering means 45 and the second decoding means 47. In the event that the session key is changed, a new random number r2 (not shown) is generated by the random number generating means 43; and in the same manner as described above, r2 is held in common by the cipher communication by means of the master key km and is made to the session key.\nNow, the case of cipher communication between two communication parties A and B is assumed. The communication party A and the communication party B encipher the session key r1 by using the master key km and send. When another communication party C holds the master key km, all cipher text which is exchanged between the communication parties A and B can be deciphered since the communication party C can decipher the session key r1. Therefore, it is required that the master key km is known by only the communication party A and the communication party B, and in the event that, for example, the communication party A carries out cipher communication with the communication party C, a master key which is other than km is used. Namely, the communication party A must hold the master keys which are identical with the number of the parties which carry out the cipher communication.\nIn such conventional system, there is no problem in the event of a small number of parties are to be communicated; but in a network having unspecified many subscribers, communication parties to be communicated become large number, and management of the key becomes a big problem. As a means for solving it, the method in which a center for performing management of the key is provided, and prior to prosecution of the cipher communication, the center delivers (or transmit) the common session key to both the parties, is generalized, but there is a defect that the center must intervene in every common holding operation of the key. On the other hand, as other solution, there is a method using the public key cipher system which is superior in management of the key, but the public key cipher system entails a much longer processing time in comparison with the secret key cipher system. As mentioned above, in the conventional cipher communication method, a big problem has existed with respect to the management of the key.\nIn view of such a problem, the present invention is directed to provide a data carrier which is safely, easily and speedy realizable common ownership of the key, on the basis of the futures in which an internal data is physically safety and calculation ability exist, and a data communication apparatus using it."}
{"text": "This invention relates to an electrically powered actuator of the type which includes a bidirectional electrically powered drive member a driven member coupled to the electrically powered drive member to be driven thereby, and a travelling member engaged with the driven member so as to be displaceable in opposite directions relative to the driven member between limits of displacement and in accordance with direction of drive of the drive member.\nA known electrically powered actuator of the type described above is manufactured by Fortune Engineering Limited of Blantyre, Glasgow G72 9BR, Scotland, and marketed under the trade mark Rolaram.\nAn application in which it is desirable to use an electrically powered actuator is in a device for supporting and moving the cabin of a vehicle simulator such as a flight simulator. However, because the total moving weight in such an application may be as much as 17,000 kilograms and may be supported by six actuators to give six degrees of freedom, it is essential that the device for supporting and moving the weight be safe in the event of an electrical supply failure or cutting off of the supply to the electrically powered drive member."}
{"text": "The invention relates to coils and has particular application to wound bobbin coil apparatus typically used in small electric motors.\nWound bobbins are widely used in the fabrication of small electric motors. The bobbins are typically automatically wound, by various types of machinery, with relatively small copper wire. Such copper wire typically is provided with only a varnish type insulation. The wires are provided with extensions or \"terminations\" which are typically multistrand wires having a substantially greater outside diameter and having conventional insulation disposed over most of the axial extent thereof. Such terminations are necessary to provide durability which would not be possible with the relatively fragile wire which is used for the windings of the coil itself.\nThe prior art techniques for providing such terminations commonly include the utilization of a so-called \"saddle\" which holds the terminations (typically two in number) while the various connections are made. While the prior art method has been widely used it does not appear to offer the maximum efficiency.\nIt is an object of the invention to provide a coil construction which will be extremely durable and which more specifically, will provide substantial strain relief for the terminations.\nIt is another object of the invention to provide a construction which will require substantially less labor than is necessary to assemble apparatus of a more conventional design.\nYet another object of the invention is to provide apparatus which can be manufactured at a relatively low cost both for labor and material."}
{"text": "The present invention relates to a technique for adding an inherent item, which is to be referred to or to be updated by a particular organization, to information (such as a table) shared among a plurality of organizations in a Relational Data Base Management System (RDBMS).\nThe RDBMS is a system to manage data shared among a plurality of organizations. For example, in an information system of a firm or a company, corporate information of firms as its customer is stored in the RDBMS to share the corporate information among organizations (e.g., divisions) of the firm. The corporate information is stored in the RDBMS in the form of a table including a plurality of items (columns). A business application (AP; program) accesses the table to refer to or to update the corporate information. The items are, for example, names and locations of the firms.\nIn such data processing system, the respective organizations desire to use mutually different items depending on cases. FIG. 2 shows an example of this situation in a customer corporate information management system. In the customer corporate information management system 100, in-house organizations such as a sales division and a supply division share customer corporate information 130. However, the respective divisions desire to use different items. A customer management application (sales division customer management AP 110) of the sales division refers to information 111 to use items which store a salesperson and a contact address, e.g., a contact phone number (salesperson contact phone number) of the salesperson of an associated firm. On the other hand, an order application (supply division order AP 120) of the supply division refers to information 121 to use a state of approval and an approval number of International Organization for Standardization (ISO) 9001, i.e., the international standard for quality certification, to thereby determine whether or not a firm under consideration is appropriate as a supplier.\nAssume that items shared among the organizations are referred to as shared items and items used only by a particular organization are referred to as inherent items. In the example of FIG. 2, the firm names and the firm locations are shared items. The salesperson and the salesperson contact phone number are inherent items used only by the sales division.\nIt is not assumed that the one-to-one correspondence exists between the business applications (sales division customer management AP 110 and supply division order AP 120 in this example) and the organizations. Depending on cases, a business application is used by only one organization or by a plurality of organizations. In the former situation, the business application uses inherent items of a particular organization. In the latter case, the business application uses inherent items of organizations to which the users of the application belong. That is, even one business application uses different inherent items depending on users.\nIn the conventional RDBMS, the inherent items are managed in three methods as below.\nAccording to a first method, for each organization, a table is prepared to store inherent items (which will be described in conjunction with FIG. 25). In a second method, for each organization, a table is prepared to store both of shared items and inherent items (which will be described in conjunction with FIG. 26). In a third method, a table is prepared to store both of shared items and inherent items for the respective organizations (which will be described in conjunction with FIG. 27).\nThese three methods have been introduced as three mapping patterns provided by Hibernate which is an object/relational mapping tool in wide use, in Section 3.6 “Class Inheritance Mapping” of “Hibernate In Action”. Reference is to be made to pages 121 to 131 of “Hibernate In Action” written by Christian Bauer et al, translated by Akira Kurahashi et al, and published from Softbank Creative on Jan. 15, 2006.\nMethods 1 to 3 respectively correspond to “3.6.3 Table per subclass”, “3.6.1 Table per concrete class”, and “3.6.2 Table per class hierarchy” of the article described above.\nIn general, according to the class structure (logical data structure) of objects to be mapped, shared items are designated as attributes of a main class and inherent items of each organization are designated as attributes of a subclass for each organization. FIG. 24 shows an example of the class structure for the data items of FIG. 2. The main class having a shared item as an attribute is associated with an in-house shared corporate attribute 2401, and the subclass having an inherent item as an attribute is associated with a sales division inherent corporate attribute 2402 and a supply division inherent corporate attribute 2403. The sales division inherent corporate attribute 2402 has, as attributes, inherent items to be used by the sales division. The supply division inherent corporate attribute 2403 has, as attributes, inherent items to be used by the supply division.\nDescription will be given of FIGS. 25 to 21 based on the above assumption. In FIG. 25, shared items are stored in an in-house shared table 2501, inherent items to be used by the sales division are stored in the sales division inherent table 2502, and inherent items to be used by the supply division are stored in the supply division inherent table 2503.\nIn FIG. 26, shared items 2602 and inherent items to be used by the sales division are stored in a sales division inherent table 2601, and shared items 2612 and inherent items to be used by the supply division are stored in the supply division inherent table 2611.\nIn FIG. 27, a shared item 2711, a sales division inherent item 2712, and a supply division inherent item 2713 are stored in one in-house shared table 2701."}
{"text": "The present invention relates to a semiconductor device comprising a chip size package of the type used for a high density mounting module, a multichip module, and the like, a method of manufacturing the same, and to a semiconductor wafer used for manufacturing the semiconductor device.\nCurrently, in accordance with the current trend toward making electronic devices small, while retaining a high performance, semiconductor devices used therein are required to have a high integration, a high density, and to have a fast processing speed.\nCorresponding to such a requirement, the methods of manufacture of semiconductor devices have been changed from the pin insertion type to a surface implementing method for increasing the implementing density; and, in order to accommodate an increasing number of pins, packages, such as a DIP (Dual Inline Package),a QFP (Quad Flat Package), a PGA (Pin Grid Array), and the like, have been developed.\nHowever, implementation of the QFP is becoming difficult because of the increasing number of pins, since the connecting leads with the implementing substrate are concentrated at the periphery of the package, and the lead itself is thin and can be readily deformed. The PGA is disadvantageous in a high density implementation, because the increasing processing velocity is difficult to attain electrically, since terminals of the PGA for connecting with the implementing substrate are slender and are concentrated so as to be extremely close to each other, and surface implementation is impossible, since the PGA package is a pin insertion type.\nRecently, in order to solve the above problems and to realize a semiconductor device having an increased processing velocity, a BGA (Ball Grid Array) package has been developed. The BGA package comprises a stress buffer layer between semiconductor chips and a substrate whereon circuits are formed, and bump electrodes, which are external terminals, on the implementing substrate side of the substrate whereon circuits are formed (U.S. Pat. No. 5,148,265). A package having a BGA structure is readily implemented on the surface, because deformation of the leads, such as occurs in a QFP does not occur, since the terminals for connection with the implementation substrate are provided in the form of ball-shaped solders, and the pitch between the terminals can be wide, since the terminals are dispersed all through the implementing plane. Furthermore, because the length of the bump electrode, i.e. the external terminal, is short in comparison with the PGA package, the inductance component is small, the signal velocity is fast, and it is possible to accommodate the requirement to increase the processing velocity.\nOn the other hand, JP-A-8-172,159 (1996) discloses a LOC (Lead On Chip) package, which comprises a cross sectional composition formed of sealing material/chip/protecting film/sealing material, as a chip provided with a protecting film. The protecting film increases the adhesion of the chip with the sealing material, and, at the same time, protects the chip from damage by a pick up pin.\nJP-A-7-135189 (1995) discloses an invention relating to a wafer adhesion sheet for manufacturing semiconductor devices having a LOC structure. The wafer adhesion sheet is used as a protecting film until the chip is mounted in a package during the process for manufacturing the semiconductor.\nRecently, in accordance with widespread use of portable information terminals, a decrease in size and a high density mounting of semiconductor devices are required. Therefore, a CSP (Chip Scale Package), the package size of which is almost the same as chip, has been developed. A CSP of various types has been disclosed in “Nikkei Micro Device” p38–p64, published by Nikkei BP Co. (February, 1998). A CSP is typically manufactured by the steps of: adhering semiconductor chips, which has been cut into respective pieces, onto a polyimide substrate or ceramic substrate, whereon a circuit layer is formed; electrically connecting the circuit layer with the semiconductor chip by a method such as wire bonding, single point bonding, gag bonding, bump bonding, and the like; sealing the connecting portion with a resin; and forming external terminals, such as solder bumps.\nJP-A-9-232256 (1997) and JP-A-10-27827 (1998) disclose methods of mass production of a CSP. In accordance with these methods, the semiconductor device is manufactured by the steps of: forming bumps on the semiconductor wafer; connecting the circuit substrate electrically via the bumps; sealing the connecting portions with a resin; forming the external electrodes on the circuit substrate; and cutting the wafer into respective pieces. The “Nikkei Micro Device” p164–p167, published by Nikkei BP Co. (April, 1998), discloses another mass production method for the CSP. In accordance with the disclosed method, the semiconductor device is manufactured by the steps of: forming bumps on the semiconductor wafer by soldering; sealing portions other than the bumps with a resin; forming external electrodes at the bump portions; and cutting the wafer into respective pieces.\nAmong the CSP which are assembled by adhering the semiconductor chips cut in pieces to a polyimide substrate or a ceramic substrate, a CSP wherein a circuit layer is connected to the chip by wire bonding becomes larger than the chip size, naturally, because the bonding area of the circuit layer is located outside of the chip. In the case of a CSP connected by bump bonding, the substrate becomes larger than the chip in order to prevent resin from flowing down at the time of potting, because the interval between the chip and the substrate is sealed with resin after bonding. Accordingly, a problem occurred in that the package size of such a CSP became larger than the chip.\nA CSP using chips cut in pieces experienced a problem in that it took a long time to manufacture the semiconductor device, because, after dicing the chips, each respective chip must be located correctly on the substrate, adhered thereon, connected electrically, and sealed.\nA CSP using a resin substrate, such as a substrate made of polyimide and the like, as the circuit layer had a problem in that water, absorbed in the package at re-flowing when the package was provided onto the mounting substrate, was expanded and failures, such as bubble formation and peeling off, were generated, because the chip was adhered with an adhering agent.\nA CSP, which was manufactured by the steps of: forming bumps on the semiconductor wafer; connecting the semiconductor wafer with the substrate; sealing the interval between the substrate and the semiconductor wafer with a resin; forming the external electrodes; and cutting the semiconductor wafer into respective pieces; had a problem of warping of the semiconductor wafer and the semiconductor device due to curing shrinkage, because the resin layer was formed on only one side of the wafer.\nAdditionally, except for the wire bonding type CSP, many of the CSP have an exposed plane, which is opposite to the plane whereon the circuits are formed on the chip. Therefore, there was a problem in that failures, such as cracks, were generated at the edge of the chip, and damage to the rear plane occurred due to the package falling down during transportation and handling thereof, such as when the package is picked up during a mounting operation."}
{"text": "At the present time several magazines run a dating column where a person may leave a message for an unknown person whom they have seen or encountered, in the hope that said person will see the message and respond to them. This system is unsatisfactory because of the low probability that the person to whom they wish to send the message would see the message, the low probability that they would be able to recognise that it was themselves for whom it was intended and the fact that because they would not see it until some time after the chance meeting, they are less likely to want to respond.\nThe aim of the present invention is to allow a person to send a message to a stranger whom they encounter. Typically, they will send an e-mail or short text message from their own mobile telephone.\nThe invention aims to enable that message to reach the intended person and preferably to allow them to respond in a fun, safe and convenient fashion."}
{"text": "1. Field of the Invention\nThe present invention relates generally to traveling lights, and more particularly to lights for use with in-line roller skates.\n2. Description of the Prior Art\nRoller skating at night or near dusk is an exciting but potentially dangerous activity. Recognizing this fact, inventors have tried to provide lights for roller skates, but have failed to provide lights that are genuinely easy to apply to, and remove from, a variety of skates.\nU.S. Pat. No. 4,367,515, issued on Jan. 4, 1983, to Steven F. Beard, describes a roller skate light attachment that mounts on a stop portion of a roller skate. There is no provision for mounting the attachment on a skate not having a stop portion.\nU.S. Pat. No. 4,463,412, issued on Jul. 31, 1984, to Ronald W. Broach, describes an illuminated shoe skate attachment. The attachment mounts to bolts or ridges on a skate by use of a complicated set of brackets. The attachment is detachable for use on other skates, but only if the brackets fit the other skates.\nU.S. Pat. No. 4,991,066, issued on Feb. 5, 1991, to Gregory L. McCowan, describes an adaptor kit for providing lighting for a skateboard. The lighting replaces a part ordinarily present on a skateboard and is thus unsuitable for use with roller skates.\nU.S. Pat. No. 4,997,196, issued on Mar. 5, 1991, to John L. Wood, describes a system in which lights are mounted in grooves in a skateboard or in skid bars in a skateboard. There is no provision for mounting lights on roller skates.\nU.S. Pat. No. 5,033,212, issued on Jul. 23, 1991, to Walter R. Evanyk, describes a lighting system which is mounted to objects such as shoes by use of hook-and-loop fasteners. There is no use of stretchable cord to ensure an easy and lasting engagement of the lighting system with an object to which it is attached.\nU.S. Pat. No. 5,327,329, issued on Jul. 5, 1994, to David L. Stiles, describes lighting attachments for in-line skates that are frictionally and adhesively mounted. There is no use of stretchable cord to ensure an easy and lasting engagement of the lighting system with an object to which it is attached.\nCanadian Patent Document No. 1 253 832, issued on May 9, 1989, to Nicholas A. Rodgers, describes lights for footwear. There is an internally disposed switch that responds to movement of a wearer. There is no provision for mounting lights on roller skates.\nFrench Patent Document No. 2 643 794, issued on Sep. 7, 1990, to Jean Leonard Darfeuille, describes head-lights for footwear. There is an internally disposed switch that senses presence of a wearer. There is no provision for mounting lights on roller skates.\nNone of the above inventions and patents, taken either singly or in combination, is seen to describe the instant invention as claimed."}
{"text": "The invention relates to a trampoline with a base, which includes a frame, and with a jumping sheet, which preferably includes hook elements on its underside, the jumping sheet being connected to the frame by means of a plurality of annular strap sections, which each extend at least partially around the frame and are secured to the hook elements. The invention is also directed to a hook element for trampolines.\nTrampolines are widely used in physiotherapy and in the leisure area.\nIn the case of so-called mini-trampolines, a distinction is made between trampolines with a spring suspension and a rubber strap ring suspension. In the case of a spring suspension, the jumping sheet is secured to the frame by means of spiral springs. Such a suspension makes considerable acceleration forces available and is thus suitable, in particular, for sports training with high demands on the cardiovascular system.\nIn the case of a rubber strap ring suspension, one or more elastic straps or strap sections are used, which connect the jumping sheet to the frame. The larger elasticity of such a fastening produces “softer” deceleration of the body and is thus primarily suitable in the case of joint and back problems, in rehabilitation, for relaxation and combating stress and for children.\nOnly a single elastic strap was originally provided in strap suspensions, see for example DE 299 19912 U1. The strap is wound alternately about the frame of the trampoline and pulled through loops arranged on the underside of the jumping sheet. It has, however, transpired in practice that the known trampoline has disadvantages in operation. Thus the fastening of the jumping sheet to the frame is perceived to be complex. Furthermore, tilting of the jumping sheet with respect to the frame can occur in use. In particular, however, the entire strap must be replaced if only one point on the strap is abraded or worn.\nThere has therefore been a move to suspending the jumping sheet on a plurality of annular strap sections, which are preferably fastened to the underside of the jumping sheet and are looped around the frame. Such a trampoline is disclosed in, for instance, DE 102 26707 B4 emanating from the applicant. By contrast with the suspension of the jumping sheet with only a single strap, the individual strap suspension combines a number of advantages. Above all, the durability is significantly increased. Furthermore, replaceability is also better (and cheaper). If a strap section is worn, it is simply replaced individually. The other strap sections remain untouched. Finally, the individual strap suspension always ensures correct centering of the jumping sheet.\nThe individual strap suspension has basically proved to be satisfactory. There is, however, the endeavour to improve the known trampoline. This applies, in particular, with regard to the wear resistance of the suspension."}
{"text": "1. Technical Field\nThe present invention is directed to elevator dispatching. More particularly, the present invention is directed to assigning an elevator car to a floor in response to a hall call registered at the floor, based on a series of bonuses and penalties and remaining response time.\nAs used herein, remaining response time means an estimation of the amount of time required for an elevator car to reach the floor at which the hall call is registered, given the car calls and hall calls to which the elevator car is committed.\n2. Background Information\nIn a building having a plurality of floors, each floor typically has a set of buttons located in the hallway at or near the elevators. These buttons, commonly referred to as hall call buttons, enable users to request elevator car service in a predetermined direction, i.e., up and/or down. Additionally, the interior of an elevator car is generally equipped with a plurality of buttons, commonly referred to as car call buttons, which enable users to request service to specific floors.\nIn simplified terms, an elevator control system, also referred to in the art as an elevator dispatching system, monitors the status of the hall call buttons at the floors and car call buttons in the elevator cars, assigning elevator cars to the floors in response to hall call and/or car call registration.\nAs used herein, relative system response, commonly abbreviated RSR, means an elevator dispatching system which employs a series of bonuses and/or penalties to determine which elevator car to assign to a registered hall call. RSR dispatching systems are well known in the art. For example, U.S. Pat. No. 4,363,381 issued to Bittar, entitled Relative System Response Elevator Call Assignments, U.S. Pat. No. 4,815,568 to Bittar, entitled Weighted Relative System Response Elevator Car Assignment System With Variable Bonuses And Penalties, U.S. Pat. No. 4,782,921 to MacDonald et al., entitled Coincident Call Optimization In An Elevator Dispatching System, U.S. Pat. No. 4,790,412 to MacDonald et al., entitled Anti-Bunching Method For Dispatching Elevator Cars, and U.S. Pat. No. 4,793,443 to MacDonald et al., entitled Dynamic Assignment Switching In The Dispatching Of Elevator Cars, herein incorporated by reference.\nIn RSR dispatching systems, each elevator car has associated therewith a value for each one of the bonuses and penalties, the values being dependent upon the elevator status, relative to the registered hall call. An RSR value for each elevator car is determined by cumulating the bonuses and penalties for each elevator car, and the elevator car having the most favorable RSR value is assigned to respond to the registered hall call.\nCertain penalties of the Bittar systems combine to yield an estimate of the response time of each elevator car to respond to the registered hall call. Some of the various penalties used in the Bittar systems include a run time penalty (RTP) based on the time required for an elevator car to travel between hall and car stops assigned thereto; a travel-through express zone penalty (TRE) based on the time required to travel through the express zone; and a hall stop penalty (HSP) and a car stop penalty (CSP) based on the time delay incurred for each hall stop and car stop, respectively, assigned to the elevator car.\nThe bonuses of the Bittar systems offset the penalties, thereby favoring an elevator car based on certain conditions. Examples of various bonuses used in the Bittar systems include a coincident car call bonus (CCB) which favors an elevator car having a car call coincident with the floor registering the hall call; a contiguous stop bonus (CSB) which favors an elevator car having a commitment at a floor contiguous with the floor registering the hall call; and a previously-assigned bonus (PAB) which favors the elevator car previously assigned to the hall call.\nIn the Bittar systems, the RSR value is determined by subtracting the bonuses and adding the penalties, and the elevator car having the lowest RSR value is the most favorable. Similarly, in the MacDonald systems, the RSR value is determined by substracting bonuses from an estimated arrival time, based on items substantially similar to RTP, TRE, HSP and CSP, with the elevator car having the lowest RSR value being the most favorable. The bonuses employed by the MacDonald systems are similar to the CCB (see U.S. Pat. No. 4,782,921), the CSB (see U.S. Pat. No. 4,790,412), and the PAB (see U.S. Pat. No. 4,793,443) of the Bittar systems.\nAlthough these systems are defined such that the lowest RSR value is most favorable, additions and subtractions can be reversed such that the highest RSR value is most favorable.\nThe RSR determination typically occurs either every cycle, e.g., every 250 milliseconds, or on an \"as needed\" basis, e.g., whenever an elevator car changes positions, responds to a hall or car call, or whenever a new hall or car call is registered. It is conceivable that the system can reassign the unanswered hall calls to a different elevator car quite often. In order to dampen what might otherwise be an erratic dispatching system, the previously assigned bonus (PAB) is included in the Bittar and MacDonald RSR systems to favor the elevator car which was previously assigned to the unanswered hall call.\nThe habits and customs of the elevator user dictate the value typically assigned to the previously assigned bonus (PAB). The user predominantly in Europe and the Americas wants to be informed which elevator car will be arriving in response to the hall call shortly before it arrives. Thus, the hall lantern located at or near the arriving elevator car illuminates and/or sounds shortly before the elevator car arrives. In these types of RSR systems, the Bittar systems typically assign a relatively low value to the PAB. Thus, if another elevator car can service an unanswered hall call by at least the PAB value before the assigned car, the unanswered hall call is reassigned to that other car.\nThe user predominantly in Japan, on the other hand, wants to be informed which elevator car will be responding to the hall call at the time of hall call registration. In this way, the user can wait at the door of the assigned elevator. In these RSR systems, commonly referred to as instantaneous car assignment systems, the elevator car having the most favorable RSR value is almost immediately assigned to the hall call, and the hall lantern located at or near the assigned elevator car illuminates and/or sounds upon assignment. In these types of RSR systems, the Bittar systems typically assign a relatively high value to the PAB, so as to maintain the integrity of the initial car assignment.\nThere are situations where the initial elevator car assignment, although the best when made, subsequently turns out to be less than satisfactory. For example, in response to an intervening car and/or hall call, the boarding passengers register a plurality of car calls. Additionally, the elevator car, while traveling through the lobby or while responding to an intervening car and/or hall call, gets delayed by boarding passengers who hold the door open, e.g., to wait for others or to finish a conversation with a non-boarding person. Further, an empty elevator car may be assigned a hall call moments before a boarding passenger registers a car call in the direction opposite the assigned hall call.\nDue to the large value of PAB, in immediate car assignment systems any one of these situations can drastically extend the system's registration time, defined as the time between when the hall call is registered and when the elevator car is about to arrive in response thereto, as indicated by the hall lantern at or near the arriving elevator car. The maximum registration time is a conventional indicium of overall system responsiveness. In RSR systems employing instantaneous car assignment, some buildings have rather large maximum registration times.\nSuch large registration times are considered highly unacceptable for at least two reasons. First, a user's irritation level is a function of the amount of time spent waiting for an elevator car. Thus, the longer the wait, the more severe his or her irritation. Second, and more importantly, before the assigned elevator car can get to the floor to service the unanswered hall call, the floor may be bypassed, up-going hall call, by at least one of the other elevator cars; e.g., a car traveling in the up direction travels past a floor having an unanswered up-going hall call.\nHowever, due to the relatively large PAB value, the initially-assigned elevator car remains assigned to the unanswered hall call, despite the fact that at least one other elevator car is bypassing the floor having the unanswered hall call. The other elevator cars do not stop because they do not get the assignment."}
{"text": "The present invention relates to polymer solutions containing a sulfur-containing or oxygen-containing aromatic polymer, methods of forming such solutions, methods of using such solutions to form conducting polymer articles including films and methods of using such solutions as electrically conducting liquids.\nConductive poly(p-phenylene sulfide) and other chalcogenide polymers rendered conductive by doping with a Lewis Acid halide and other dopants are disclosed in European published patent application No. 31,444 (July 8, 1981) the subject matter of which is included in U.S. Pat. No. 4,375,427, issued Mar. 1, 1983. Such materials are prepared, in some forms, by introducing a gaseous dopant into the solid polymer pellet.\nIt is indicated in U.S. Pat. No. 4,375,427 that arsenic trifluoride gas, when introduced with arsenic pentafluoride gas as dopant, increases the doping efficiency.\nIn the general field of conducting polymers, such as polyphenylene, polyphenylene sulfide, polypyrrole and polyacetylene, it has been heretofore believed impossible to dope one of these polymers to the extent that it becomes a semiconductor or conductor and thereafter dissolve the polymer. In the absence of such solutions, the ability to form certain articles out of conducting polymers, and especially conducting or semiconducting polymer films, is restricted. In particular, doping with a gas after forming a polymer article is known to expand or distort the polymer article. Thus a need exists for techniques to facilitate the production of conductive polymers, to facilitate the doping process and, especially, to form shaped articles such as films and fibers of conducting polymers."}
{"text": "Sweet syrups are widely used in the baking, confectionery and beverage industries, for example. These syrups generally consist of sucrose (cane sugar) or dextrose-containing products obtained from starch hydrolysis as the principal sweetening agent. When a syrup is needed that is sweeter than that obtained from sucrose, an invert sugar syrup is employed. This is produced by acid hydrolysis of sucrose to produce a mixture of about 50 percent glucose (dextrose) and about 50 percent fructose (levulose). While glucose is somewhat less sweet than sucrose, the fructose is considerably sweeter than sucrose so that the overall sweetness is increased as compared to sucrose.\nIt is well known that dextrose can be converted under alkaline conditions to fructose. This conversion has great potential value in the production of sweet syrups. However, the alkaline conversion has not been commercially successful because the alkaline reaction produces an undesirably high ash level in the product syrup which is uneconomical to remove. The syrup is not acceptable unless this ash is removed. Alkaline isomerization and its attendent problems are discussed in U.S. Pat. No. 3,383,245.\nThe prior art then turned to an enzyme conversion of glucose to fructose. It was found that species of Pseudomonas hydrophila, Streptomyces flavovirens, Streptomyces achromogenus, Streptomyces echinatur, Streptomyces albus, Streptomyces olivaceus, and Bacillus coagulans, for example, could be grown in appropriate nutrient media to form enzymes having glucose isomerase properties. This is described in U.S. Pat. Nos. 2,950,228; 3,616,221; 3,625,828; and 3,979,261, for example. None of the known prior art suggests the use of a Bacillus licheniformis species to produce a glucose isomerase having the desirable characteristics of being stable under a wide range of pH and temperature conditions."}
{"text": "This invention relates to a process for producing a single crystal of garnet ferrites, and more specifically to a process for producing a single crystal of garnet ferrite suitable for use as a Faraday rotation material in an optical isolator adapted to prevent light from returning to a semiconductor laser employed for optical communication or in an optical magnetic field sensor adapted to control or protect electric power system in transformer substations.\nAs techniques for producing a single crystal of garnet ferrite useful as magnetooptical materials, the flux technique and LPE technique have been usually employed to date.\nOf these techniques, the flux technique comprises melting raw materials for such respective components that form a ferrite having the garnet structure together with a flux component such as PbO, PbF.sub.2 and/or B.sub.2 O.sub.3 in a crucible, for example, made of platinum into a homogeneous melt and cooling the melt gradually to cause a single crystal of the garnet ferrite to grow.\nFurther, the LPE technique also comprises similarly to the flux technique, melting raw materials for such respective components that form a ferrite having the garnet structure into a homogeneous melt as described above and dipping a substrate having the garnet structure in the resulting melt while maintaining it at a temperature slightly lower than the saturation temperature to cause a single crystal of the garnet ferrite to grow on the substrate as a thin film.\nHowever, the flux technique is accompanied by such problems that it requires long time, i.e., from a week to a month, for the growth of a single crystal and it cannot avoid the inclusion of impurities from the flux and crucible used.\nOn the other hand, the LPE technique requires shorter time for the growth of each crystal. However, in this technique, a larger single crystal having a three-dimensional structure cannot be obtained and mixing of impurities is also unavoidable because a flux and a crusible have to be used similarly to in the former technique.\nAmong these impurities, mixing of impurity ions other than trivalent ions, especially, Pb.sup.2+, Pt.sup.4+ and the like in a single crystal of garnet ferrite leads to a modification to the ionic valence of ferric ions (trivalent ions), thereby adversely affecting primarily the light absorption characteristic of the single crystal.\nIn addition, the above-employed crucible is usually made of a noble metal such as platinum. In order to cause a large single crystal to grow, it is indispensable to enlarge the crucible, resulting in a problem that the production cost jumps up.\nMoreover, although it has been known that the degree of Faraday rotation .theta..sub.F (deg/cm) of the garnet ferrite can be remarkably increased by substituting the c-site component in the garnet structure partially with Bi.sup.3+, and increase in Faraday rotation is useful for miniaturing an element, it has been difficult to prepare single crystals of the garnet ferrite having such a large Faraday rotation."}
{"text": "Membrane processes have been applied to a wide variety of industrial separation processes for process schemes requiring either concentration or purification of aqueous streams ranging from clarification of juices, removal of products from fermentation broths and desalination of brackish waters, etc. Some of the processes described above have been demonstrated on an industrial scale. This technology has been possible due to the development of new membranes which exhibit a high degree of hydraulic permeability coupled with ability to retain small molecules.\nThe membrane filtration process is, however, adversely affected by two phenomena, namely, concentration polarization and fouling. Concentration polarization is caused by the accumulation of solute molecules at the upstream surface of the membrane. This phenomena causes a reduction in the efficiency and of rate of the membrane filtration process. Concentration polarization is generally a hydrodynamic and diffusion phenomena.\nFouling also occurs with a variety of feeds including proteins and colloids. Fouling can be the result of an insoluble precipitate/ particle or gel layer within the pores (plugging) or the result from build-up on the surface of the membrane itself. It has been reported in the literature that up to 100 micron thick gel layers can be formed over a period of 24 hours as a result of the denaturation of proteins. The presence of fat globules can also be responsible for formation of gel layers.\nSuch gel layers and plugging are responsible for decreases in membrane flux and, hence, affect product throughput and economics. Various methods of treating fouling and concentration polarization have been reported in the literature but have only had limited success in minimizing the problem. Some of these methods are heat treatment with pH adjustment, immobilized protease treatment and membrane scouring.\nSome typical industrial applications of membrane separation in the food and beverage industry are removal of casein, fats and lactose from whey; clarification of liqueur and vodkas; sterilization of liquids, e.g. beer, wine; continuous microfiltration of vinegar; and concentration and demineralization of cheese, whey, soy whey and vegetable extracts. Other applications in the wastewater treatment industry include removal of cyanides from electroplating wastewaters; reuse of wastewater from ammunition manufacture; recovery and recycling of sewage effluent; and recovery of starch and proteins from starch factory waste, wood pulp and paper processing. Further applications are listed in an article by Minh S. Le et al, The Chemical Engineer, pp. 48-53, July/August 1985.\nPrevious work in this area includes the following U.S. Pat. Nos.: 2,571,247; 3,980,541; 4,032,454; 4,057479; 4,100,068; 4,224,135; 4,269,681; 4,276,146; 4,343,690; 4,421,579; and 4,617,128.\nMembrane separation techniques typical of the prior art include those that employ classical crossflow electrofiltration. Several good review articles that give an excellent background discussion include: Electrokinetic Membrane Processes, Milan Bier, Membrane Processes in Industry and Biomedicine, Milan Bier--Editor. Plenum Press, N.Y.-London, (1971) pp. 233-266; A Solid/Liquid Separation Process Based on Cross Flow and Electrofiltration, J. D. Henry, Jr., et al, AlChE Journal, Vol. 23, No. 6. November 1977, pp. 851-859; Membrane Fouling Prevention in Crossflow Microfiltration by the Use of Electric Fields, R. J. Wakeman et al, Chemical Engineering Science, Vol. 42, No. 4, pp. 829-842, 1987. In classical crossflow electrofiltration the material to be filtered passes between a first electrode and a membrane or filter with a second electrode positioned on the other side of the filter. The electrodes are energized to pull suspended particles from the material to be filtered toward the first electrode and away from the filter. Thus the suspended particles cannot deposit on the filter or membrane.\nU.S. Pat. No. 4,604,174 to Bollinger et al teaches an improvement of the prior art crossflow: however, the general design follows FIGS. 1 and 2 herein that illustrate typical prior art crossflow arrangements. In Bollinger, one electrode assembly forms a wall of the crossflow chamber and the second electrode is located in the filtrate chamber. This arrangement causes the slurry material passing through the crossflow chamber to flow between the two electrodes see FIG. 1 of Bollinger and the discussion in the Summary of the Invention, Column 2, line 57 to Column 3, line 32.\nObjects of the present invention include the enhancement of crossflow and dead-end electrofiltration by reducing energy requirements, improving filtration efficiency, reducing filtration time, reducing costly filter cleaning steps, and reducing filter fouling."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of vacuum cleaners in general, and in particular to an adjustable handle construction for a vacuum cleaner.\n2. Description of Related Art\nAs can be seen by reference to the following U.S. Pat. Nos. 4,662,026; 5,016,315; 5,109,568; and 5,564,160, the prior art is replete with myriad and diverse handle constructions for vacuum cleaners.\nWhile all of the aforementioned prior art constructions are more than adequate for the basic purpose and function for which they have been specifically designed, they are uniformly deficient with respect to their failure to provide a simple, efficient, and practical handle construction for vacuum cleaners that can be angularly adjustable, not only to suit the personal preferences of the individual user, but also to reduce the occurrence of carpel tunnel syndrome among individuals who use a vacuum cleaner on a daily basis.\nAs most individuals in the building maintenance field are all too well aware, even the most ergonomically designed rigid vacuum handle construction can lead to a repetitive motion injury with continued usage.\nAs a consequence of the foregoing situation, there has existed a longstanding need for a new and improved type of adjustable vacuum handle construction that cannot only be adjusted to a preferred position, but whose position can be varied to reduce the occurrence of repetitive motion injuries, and the provision of such a construction is a stated objective of the present invention."}
{"text": "A polymer electrolyte fuel cell (PEFC) can operate in a low-temperature region and has high energy conversion efficiency, and a time period required for its startup is short. In addition, the system of the PEFC can be made small and lightweight. Accordingly, the PEFC has been expected to find applications in power sources for electric vehicles, portable power sources, and household co-generation systems.\nHowever, large amounts of platinum catalysts are used in the PEFC. The use of the platinum catalysts causes an increase in cost, which is one factor that may inhibit the widespread use of the PEFC. In addition, a concern has been raised in that restriction is imposed on the PEFC in terms of platinum reserves.\nIn view of the foregoing, the development of a novel catalyst that can replace the platinum catalyst has been advanced. That is, for example, a carbon catalyst obtained by imparting a catalytic activity to a carbon material itself has been proposed (see, for example, JP 2007-026746 A, JP 2007-207662 A and JP 2008-282725 A)."}
{"text": "Polyolefin polymers are very well known in the art. Methods of manufacturing and/or processing polyolefins are likewise well known in the art. These methods of processing include but are not limited to blow molding, injection molding, pultrusion and pulforming, and extrusion. The problems associated with each of these processes are known, and seeking solutions to these problems, and/or improvements in processing conditions and equipment remains of interest to the art and industry.\nFor example, it is known that olefinic homopolymers such as high-density polyethylenes (HDPE) and polypropylenes generally have poor processability and poor impact properties. Because of the wide commercial uses of these polymers in various applications such as automotives, households, and. packaging, it is always desirable to discover new methods, additives, and/or equipment which enable these polymers to be better and more economically processed, and to improve the mechanical and other properties of these polymers and/or products made from the polymers."}
{"text": "At present, 3GPP (3rd Generation Partnership Project) has been examining LTE Advanced (hereinafter referred to as “LTE-A”) for a further increase in communication speed. SC-FDMA (Single Carrier Frequency Division Multiple Access) system with excellent PAPR (Peak to Average Power Ratio) characteristics is adopted in uplink in LTE. In LTE uplink, spatial multiplexing using MIMO (Multiple Input Multiple Output) can be used. In particular, MU-MIMO (Multi-User MIMO) that performs spatial multiplexing in which multiple mobile station apparatuses use the same frequency and time resource can be used.\nTo reduce degradation of characteristics in MIMO communication caused by noise, interference, or the like, it is desirable to highly accurately calculate the channel (may also be referred to as “propagation path”) of each antenna port (defined for every physical antenna or combination of two or more physical antennas) in transmission and reception. In LTE, orthogonality in frequency domain is realized by code multiplexing implemented by allocating, to each antenna port, a sequence generated by applying a cyclic shift in time domain to an uplink demodulation reference signal (UL DMRS: Uplink Demodulation Reference Signal) generated on the basis of a Zadoff-Chu sequence. Further, since there are two sequences of UL DMRS (that is, two SC-FDMA symbols) in one subframe, LTE-A applies a Walsh sequence in units of SC-FDMA symbols, thereby realizing orthogonality in time domain (Non Patent Literature 1). This spread sequence applied in time domain is referred to as an orthogonal cover code or an OCC (Orthogonal Cover Code).\nA condition in which orthogonality is established by a cyclic shift in UL DMRS is limited to the case in which the allocated frequency regions are completely the same. Therefore, when MU-MIMO is applied to mobile station apparatuses to which different frequency regions are allocated, orthogonality between the mobile station apparatuses can be realized by allocating orthogonal OCCs, such as [+1 +1] and [−1 −1], to the mobile station apparatuses, respectively. Further, as in Non Patent Literature 2, it has been proposed to use IFDM (Interleaved Frequency Domain Multiplexing) that arranges DMRS in a comb-shape on the frequency axis. Orthogonality of multiplexed UL DMRS can be realized by changing the offset of a subcarrier to be allocated."}
{"text": "Bone morphogenetic proteins (BMPs) are members of the TGF-β superfamily that can induce endochondral bone formation in adult animals. This superfamily includes a large group of structurally related signaling proteins that are secreted as dimers and then cleaved to result in biologically active carboxy terminal domains of the proteins. These bioactive proteins are characterized by 7 highly conserved cysteine residues. Interestingly, these proteins have different roles at various stages of embryogenesis and in adult animals. Recombinant BMPs are now available and have been shown to induce endochondral bone formation when assayed in vivo.\nIndeed, the initial discovery of the BMPs was facilitated by such in vivo assays for cartilage and bone development. These assays were based on the observation that bone development could be initiated by subcutaneous or intramuscular implantation of compositions comprising an extract of demineralized bone and residual bone powder. The novel proteins identified in the extracts were termed “bone morphogenetic proteins.” These proteins were subsequently classified as members of the TGF-β superfamily by virtue of amino acid sequence relatedness. Screening of genomic and cDNA libraries led to the isolation of polynucleotides encoding BMP-2, -3, -4, -5, -6 and -7.\nOne deficiency of the bone induction assay regards its inability to distinguish the physiological roles of different BMP family members. The cartilage and bone inducing activity of the BMPs is remarkable because the normal stages of endochondral bone formation that occur during ontogeny are recapitulated in the adult animal. These stages include mesenchymal condensation, cartilage and bone and bone marrow formation and eventual mineralization to produce mature bone.\nSeveral observations suggest that BMPs have wide-ranging extraskeletal roles in development. First, localization studies in both human and mouse tissues have demonstrated high levels of mRNA expression and protein synthesis for various BMPs in kidney (BMPs-3, -4, -7), lung (BMPs-3, -4, -5, -6), small intestine (BMPs-3, -4, -7), heart (BMPs-2, -4, -6, (BMPs-2, -4, -5, -7) and teeth (BMPs-3, -4, -7). Second, several members of the family, including BMP-4 and -7, are key molecules in epithelial-mesenchymal interactions, for instance during odontogenesis. Third, BMP-2 and BMP-4 are involved in the signaling pathway that controls patterning in the developing chick limb and BMP-4 is a ventralizing factor in early Xenopus development. Fourth. Drosophila homologs of the BMPs, the decapentaplegic (dpp) and 60 A gene products, have the capacity to induce bone in mammals whereas human BMP-4 confers normal embryonic dorso-ventral patterning in Drosophila transformants defective in dpp expression. Thus, the BMPs are now appreciated as pleiotropic cytokines.\nInterestingly, none of the known BMPs are strongly expressed in the chondroblasts and chondrocytes of the cartilage core of developing long bones. The hypertrophic chondrocytes, where both Vgr-1 (BMP-6, (Lyons et al., Development 109:833 (1990)) and OP-1 (BMP-7)(Vukicevic et al., Biochem. Biophys. Res. Commun. 198:693 (1994)) have been found are exceptions in this regard."}
{"text": "The present invention is related to transportable service carts, of the type commonly used for the convenient storage and transportation of items.\nService carts of the described type come in a variety of shapes and sizes that are typically tailored to a specific application or for use with particular items. They usually comprise a storage compartment and a plurality of wheels. The storage compartment may also comprise surfaces and/or shelves that provide a support surface for the stored items and may comprise drawers, cabinets, or other components that may completely enclose the stored items. The plurality of wheels is provided for convenient transport of the service cart.\nOne application for a service cart is in automotive repair facilities, wherein the service cart stores a technician's tools. The service cart may be transported to the approximate area that the technician is working, typically near a vehicle, to provide convenient access to the stored tools. Occasionally, a service cart will strike a vehicle while the vehicle is in the automotive repair facility, which is undesirable because such impacts between the service cart and vehicle may damage the exterior of the vehicle or the service cart and may damage the internal components of the vehicle or the stored items in the service cart. Repair facility equipment may also come into contact with the service cart on occasion. Such equipment may comprise shelving, vehicle lift systems, or other machinery that may also be damaged if impacted by a service cart.\nThe damage to the struck object or the service cart is typically caused by the extremities of the cart, such as the corners of the cart, protruding shelves, protruding handles, or the like. The extremities of the cart are most likely to come into contact with other objects, and the extremities often define relatively sharp features, such as a non-radiused corner or a corner with minimal curvature. Such sharp features concentrate the impact on the struck object during the impact, which increases the amount of damage to the object rather than dispersing the impact over a larger area and reducing the damage.\nTherefore a need exists for a transportable service cart that reduces or prevents damage to struck objects and the service cart when the service cart impacts other objects. A further need exists for a transportable service cart that reduces the likelihood of impacts between a service cart and other objects."}
{"text": "1. Field of the Invention\nThis invention relates to ion implantation and more particularly to a non-mass-analyzed ion implantation process wherein ions of non-selected species are accelerated and implanted into a target.\n2. Description of Prior Art\nToday, ion implantation is evaluated as one of important techniques, especially, in the semiconductor industry. Conventionally, most ion implantation apparatus for processing semiconductors have been based on mass-analyzed ion implantation wherein mass analysis is effected by utilizing a static magnetic field generated by a magnet and an ion of one selected species is used as an ion source. The mass-analyzed ion implantation has an advantage of high purity of implanted impurities but it requires a mass analyzer and suffers from a disadvantage that the apparatus is costly and running cost is high.\nOn the other hand, non-mass-analyzed ion implantation has been proposed especially for formation of a pn junction of solar batteries and it eliminates the above disadvantage (Japanese Patent Application Kokai (Laid-Open) No. 130975/78). In the non-mass-analyzed ion implantation, ions are accelerated and without being subjected to mass analysis, implanted into a target. This ion implantation process dispenses with the mass analyzer so that not only apparatus cost and running cost can be reduced but also travel path of ion beams can be reduced to reduce loss of ions due to collision and scattering of the ion beams at the inner wall of the apparatus. The reduction in loss leads to an increase in implantation current which in turn increases implantation processing speeds.\nThe conventional non-mass-analyzed ion implantation, however, simply implants all the ions of nonselected species into the target and has difficulties with the use of compound source materials. Further, even with an elemental source material such as for example red phosphorus, a phosphorus vaporizer is required and it is necessary to prevent solidification of vaporized phosphorus by heating an atmosphere between the phosphorus vaporizer and the ion source up to about 400.degree. C."}
{"text": "There are numerous reasons to limit or otherwise block the origination of outgoing calls. Parents may want to limit their children's ability to call certain people, or make certain types of calls. Further, parents may want to limit the times at which calls can be made. Office administrators or business owners often desire to impose similar limitations on employees.\nUnfortunately, there are very few ways to limit outgoing calls. Typically, a telephony switch servicing a telephony terminal must be provisioned to block certain types of calls from that telephony terminal. Such provisioning is generally limited to blocking broad categories of calls, such as 900, long distance, or international calls. In addition to being limited to blocking certain categories of calls, such provisioning requires subscribers or their representatives to request that the service providers provision the telephony switch in the desired fashion. Accordingly, current techniques for blocking certain types of outgoing calls are not only inflexible, but also cumbersome. There is a need for an efficient and flexible way for subscribers or their representatives to block outgoing calls as they desire, as well as a telephony system in which such blocking of calls can be efficiently implemented."}
{"text": "a. Field of the Invention\nThe present invention is in the field of semiconductor die packaging.\nb. Description of the Related Art\nElectronic equipment that includes packaged semiconductor dies, such as portable telephones, personal computers, and digital cameras, are being developed with increased functionality. To meet this trend, companies that assemble semiconductor dies into packages are developing semiconductor packages that include a plurality of stacked semiconductor dies in a single package.\nA conventional technique for stacking plural semiconductor dies in a single package is to couple an inactive (i.e., backside) surface of a first semiconductor die to a substrate, so that the active surface of the first semiconductor die faces away from the substrate, and then to attach the inactive surface of a second semiconductor die onto the active surface of the first semiconductor die using an electrically insulating adhesive layer. The first and second semiconductor dies are each electrically coupled to the substrate by bond wires that extend from bond pads on the active surface of the respective semiconductor dies to the substrate. In order that the upper second semiconductor die not interfere with the bond wires coupled to the lower first semiconductor die, the second semiconductor die is typically smaller than the first semiconductor die so that the second semiconductor die fits fully within the opposed parallel rows of bond pads on the active surface of the first semiconductor die. In a case where a die attach paste is used to couple the first and second semiconductor dies, the size of the second semiconductor die also is limited to account for the die-attach fillet that is formed during the die-attach process for the second semiconductor die. Unfortunately, this approach of stacking a smaller second semiconductor die on a larger first semiconductor die is not feasible where it is desired to stack two same-size semiconductor dies, or to stack a larger semiconductor die on top of a smaller semiconductor die.\nAn alternative approach for die stacking is to attach a first side of a rectangular prism of silicon or of a dielectric film, called a “spacer,” to the active surface of the first semiconductor die, and to attach the inactive surface of the second semiconductor die to an opposite second side of the spacer. The spacer has an upper and lower surface area that is smaller than the area of the active surface of the lower die, and is placed so as to be entirely within the bond pads and active surface edges of the lower semiconductor die. The spacer has a thickness sufficient to space the inactive surface of the second die above the bond wires that are coupled from the bond pads of the first semiconductor die to the substrate. Unfortunately, the spacer adds additional vertical height to the package, and thus is not optimal for certain applications.\nAccordingly, a new approach to stacking a plurality of semiconductor dies is desirable."}
{"text": "The invention relates to a method of manufacturing a lithium battery comprising a stack of a negative electrode, a separator, and a positive electrode, which method comprises the steps of applying negative electrode material on a negative current collector so as to form the negative electrode, applying positive electrode material on a positive current collector so as to form the positive electrode, and arranging a separator between the negative electrode and the positive electrode, and which method comprises the following steps:\na) producing a pattern of holes in the negative electrode;\nb) producing a pattern of holes in the positive electrode;\nc) applying a polymeric material on at least one side of the stack and subjecting the stack and the polymeric material to heat and pressure, so that the polymeric material penetrates the hole, whereby which the negative electrode, the positive electrode, and the separator are stuck and pressed together.\nThe growing market for lightweight, portable cordless consumer products, such as CD-players, mobile telephones, laptop computers and video cameras, has increased the need for high-density batteries. Specifically, very thin and flexible batteries are required. If an acceptable portability is to be achieved, the batteries contained in said consumer products should provide the necessary amount of energy at the smallest possible weight and volume. Lithium is a very advantageous material for use in batteries in which a high energy density at a minimum weight is required. However, the thinner the battery, the more difficult the application of a pressure needed to maintain a sufficient contact between the respective components of the battery.\nA method of manufacturing a lithium battery according to the preamble is known from the International patent application with publication number 00/04601.\nThe battery obtained by said method has thin and flexible shape and at the same time provides a very high energy density. Moreover, the contact between the electrodes and the separator is obtained and maintained in a very efficient way. The battery can be packed in a thinwalled canister, as the wall of such canister is not needed for maintaining a sufficient pressure on the respective components of the battery. In one of the methods according to the International application 00/04601, a film of a polymeric material is applied to both sides of the stack, and said polymeric film is subjected to heat and pressure. As a result thereof, the polymeric material melts and penetrates into the holes. A battery is obtained by said method with a polymeric film at both sides of the stack as well as polymeric material in each of the holes acting as a plug or rivet and sticking to the respective layers, causing these layers to be bonded together.\nIt is an object of the invention to provide a method of manufacturing a lithium battery according to the preamble by which an even thinner battery can be created.\nTo this end, the method of manufacturing a lithium battery according to the preamble is characterized in that the polymeric material comprises bulges which are partially located in the holes of the electrode(s).\nBy providing such bulges of polymeric material at the ends of the holes in the electrode(s), enough polymeric material is provided to fill the holes and stick together the electrodes and the separator, whereas on the outer sides of the stack, in between the holes, the polymeric material may be very thin or even be absent. The battery thus obtained can provide the same amount of energy as the battery obtained according to the prior art method, but at a smaller volume. Thus, the battery obtained according to the method of the present invention has a higher capacity.\nIn a particular embodiment, the polymeric material comprises a polymer foil which is provided with bulges.\nAs was mentioned above, the polymer foil itself may be very thin, while the bulges have to comprise enough polymeric material to fill the holes in the electrodes, thereby sticking together the electrodes and the separator.\nIn a preferred embodiment, the polymeric material is placed on a carrier foil.\nPlacement of the polymeric material on a carrier foil facilitates the handling thereof. After stacking of the battery, the carrier foil may either be removed or be left in place.\nAdvantageously, the carrier foil comprises a polymeric material with a high melting point.\nIf such carrier foil is provided, the carrier foil itself will not melt during the application of heat, whereas the polymeric material will melt, thereby penetrating the holes of the electrodes.\nIn a further preferred embodiment, the polymeric material is placed in a patterned arrangement of bulges on the carrier foil.\nDuring the application of heat the bulges will melt and penetrate into the holes in the electrodes. Afterwards the carrier foil can be removed. In the battery thus obtained, the polymeric material which has penetrated the holes in the electrodes is substantially not present at the outer sides of the electrodes.\nIn another advantageous embodiment, the polymeric material is placed in a patterned arrangement of bulges between portions of a carrier.\nAlso in this case, after the melting of the polymeric material and penetration thereof into the holes, the carrier foil can be removed from inbetween the openings of the holes in the electrodes. A battery is thereby obtained in which the polymeric material which has penetrated the holes in the electrodes is substantially not present at the outer sides of the electrodes.\nIn the above methods, use is preferably made of a porous polymeric material for the polymeric material which has to enter the holes in the electrodes, e.g. porous polyethylene. Preferably, the polymeric material is elastic. The carrier foil which may be used preferably comprises a strong foil, such as Mylar(trademark) or Kapton(trademark).\nThe electrode materials can be made by mixing negative or positive active material, conductive material, and binder material, which are all in the form of powder, in a dry process, or in a wet process in which water or an organic solvent is further added.\nThe paste-like mixture obtained is then provided over the current collector, dried and compressed. For this purpose, one of the following coating methods can be generally employed: screen printing, roller coating, doctor blade coating, knife coating, extrusion coating, bar coating, dip coating and squeeze coating. The thickness of the coated layer, which is compressed after drying, generally is in the range of 1 to 1000 xcexcm.\nPressing of the stack in the above methods is accomplished by simply pressing it for a short period of time between heated metal plates at a pressure of about 5.104 Pa at about 110 to 150xc2x0 C. The operation may also be carried out using calender rollers. After heating, the stack is cooled down to room temperature. Heating and cooling down may also be carried out in a mold with a particular shape, e.g. the shape of the appliance into which the battery is to be fitted.\nIn the same way as described above, a multilayer stack of layers can be bonded together in one step, thereby obtaining a battery of increased capacity or voltage.\nThe lithium battery of the invention can be used in various (cordless) appliances, for example notebook personal computers, portable CD-players, portable telephones, paging equipment, video cameras, electric shavers, electric tools, electric vehicles, and hearing aids. The lithium battery may be used as a primary or as a secondary battery."}
{"text": "The present invention relates to aqueous solutions and dispersions of binders for preparing aqueous coating materials, and to processes for their preparation.\nThe condensation of alcohols with carboxylic acids produces per crosslinked hydroxyl or carbonylxe2x80x94group one molecule of water, H2O:\nR1xe2x80x94OH+HOOCxe2x80x94R2xe2x86x92R1xe2x80x94Oxe2x80x94COxe2x80x94R2+H2O\nalcohol+carboxylic acidxe2x86x92ester+water\nSince the ester is the reaction product which is of interest, in the prior-art preparation processes the water produced is removed from the reaction mixture as so-called xe2x80x9cwater of reactionxe2x80x9d. This water of reaction is generally contaminated with other components of the reaction and must therefore be disposed of. A procedure of this kind is necessary, for example, in the coatings industry in connection with the preparation of coating binders.\nIn the coatings industry, the preparation of polyesters and alkyd resins from polycarboxylic acids and/or polycarboxylic anhydrides and polyalcohols and, if desired, monocarboxylic acids by condensation reactions gives rise to a quantity of about 5 to 11%, based on the polymer produced, of water of reaction. The water of reaction is normally removed effectively from the reaction mixture with the assistance of entrainers (for example xylene).\nThe vapor mixture comprising water of reaction and entrainer isxe2x80x94especially where saturated polyesters are being preparedxe2x80x94passed through a column whose function is to separate volatile polymer building blocks (especially lower polyalcohols) from the vapor mixture by liquid/vapor exchange and thereby to avoid losses of these building blocks from the reaction mixture. For polyesters and alkyd resins, which contain no constituents which are volatilexe2x80x94under the reaction conditions in the reactor or with steamxe2x80x94an overflow pipe is employed.\nIn each case, however, the vapor mixture is subsequently liquefied in a condenser and is collected in a separation receiver. In the separation receiver, an aqueous phase of relatively high specific gravity separates from the organic phase of lower specific gravity. The organic phase consists almost entirely of the entrainer employed; it is separated from the aqueous phase by means of an overflow and is passed into the reactor at the head of the column (circulation technique) possibly under control and by means of the use of a pump (possibly a metering pump).\nAnother possibility is to operate without entrainer. In this case there is no need for a separation receiver but only for a single receiver into which aqueous phase is collected.\nAlthough the aqueous phase contains high proportions of water, it is always contaminated with relatively small amounts of volatile or steam-volatile building blocks (mostly lower polyols) and/or with water-soluble degradation products from the building blocks of the polyester as well. Degradation products which may be produced and may be present in the aqueous phase of the water of reaction are monoalcohols, partial aldehydes of polyols, acetals, and in special cases cyclic, oxygen-containing compounds as well.\nSeparation of the secondary constituents from the aqueous phase during the preparation of polyesters and alkyd resins, by distillation or extraction, is complex and does not lead to reusable materials. Simple disposal of the distillate, for example by way of settling basins, is a risk. Usually, therefore, despite its small content of organic constituents the water of reaction is incinerated, usually together with other residues which are richer in organic substances. This requires separation and storage of the waters of reaction and special control as they are fed to the incinerator.\nBecause of their ecological advantages, water-dilutable coating materials have increasingly been used in the coatings industry in recent years. With these coating materials it is possible largely or completely to dispense with organic solvents whose evaporation as the coating material dries leads to emissions and to the known, associated environmental problems. In the course of preparing such aqueous coating systems, the binders of the coating material must be brought into the form of an aqueous solution or dispersion. Dispersions of this kind are highly unstable. Even small quantities of impurities may destroy or severely hinder the dispersibility of the binders. For example, a solution of this kind cannot be prepared using tap water (drinking water, mains water), which contains inorganic constituents. It is therefore necessary to employ deionized or distilled water. This implies more complex preparation and therefore increased preparation costs as well.\nThe object which the present invention has set itself is to avoid the disadvantages of the prior art, described above, which occur in particular during the disposal of water of reaction and during the preparation of binder dispersions.\nThis object is achieved in accordance with the invention by an aqueous solution or dispersion of a binder for preparing aqueous coating materials, which solution or dispersion includes as at least one aqueous component the water which can be prepared by condensation of alcohols with carboxylic acids and can be separated off from the condensation mixture.\nIt has surprisingly been found that during the preparation of polyesters or alkyd resins (which are then preferably employed for water-dilutable coating PAT 95445 systems, the water of reaction is able to replace a portion of the distilled or deionized water which is otherwise used to prepare the aqueous colloidal solution or dispersion of the polyester or alkyd resin. Unlike tap water (drinking water, mains water), which contains certain proportions of inorganic constituents which impair the usefulness for preparing aqueous colloidal solutions or dispersions of coating binders, it is possible with fractions of the waters of reaction to prepare stable aqueous colloidal solutions. The use of these aqueous polyester-resin or alkyd-resin solutions or dispersions in aqueous coating systems has only an insignificant influence on their properties. This effect was completely surprising in view of the sensitivity of the binder dispersions to impurities. On the basis of conventional experience, the opposite was to be expected. However, it is found that in most cases the secondary constituents act as cosolvents. If aqueous coating systems are employed on application lines with stoving furnacesxe2x80x94which is generally the case with the industrial use of aqueous coating systemsxe2x80x94the abovementioned substances accompanying the water of reaction do not become volatile until they are in the stoving furnace, although some of them also participate in film forming (for example reaction with amine resins). Environmental pollution as a result of the constituents, therefore, is largely ruled out, and in any case the waste air from stoving furnaces includes organic cleavage products and is therefore cleaned by means of catalytically processes or filters.\nIn the text below, the term water of reaction is intended always to denote the water which is eliminated during the condensation of alcohols with carboxylic acids and which is or can be separated off from the condensation mixture.\nIn accordance with the invention, the water of reaction is preferably taken from reactions in which the alcohols are present as polyalcohols, the term polyol referring to an organic compound which bears at least 2 hydroxyl groups. Examples of suitable polyols are ethylene glycol, propanediols, butanediols, pentanediols, neopentylglycol, hexanediols, diethylene glycol, glycerol, trimethylolethane, trimethylol-propane, pentaerythritol, dipentaerythritol, neopentylglycol hydroxypivalate, 2-methyl-2-propyl-1,3-propanediol, 2,2,4-trimethyl-l,3-pentanediol and 2,2,5-trimethyl-1,6-hexanediol It is preferred to employ neopentylglycol, 1,6-hexanediol and neopentylglycol hydroxypivalate. The polyols can of course be employed in pure form or as mixtures.\nIn the case of the novel preparation of the water of reaction, polycarboxylic acids are preferably employed. These may comprise a polycarboxylic acid containing at least three carboxyl groups, or derivatives of such an acid (e.g. anhydride, ester or halide). It is likewise possible to employ a mixture of such acids and/or acid derivatives. As examples there are, trimellitic acid, trimesic acid (1,3,5-benzenetricarboxylic acid), pyromellitic acid and trimeric fatty acids. Trimellitic acid is preferably employed.\nAs carboxylic acid it is also possible to employ a polyol which has at least one carboxyl group. Dimethylolprophonic acid is preferably employed. It is likewise possible to employ any polycarboxylic acid which contains two carboxyl groups and is suitable for preparing polyesters, or to employ a reactive derivative (e.g. anhydride, ester or halide) or a mixture of such acids and/or acid derivatives. Examples of suitable acids which may be mentioned are: phthalic acid, isophthalic acid, terephthalic acid, fumaric acid, maleic acid, endomethylenetetrahydrophthalic acid, succinic acid, adipic acid, suberic acid, acelaic acid, sebacic acid and dimeric fatty acids. It is preferred to employ phthalic acid, isophthalic acid, adipic acid and dimeric fatty acids.\nAs its aqueous component the novel aqueous solution or dispersion of a binder may include a mixture of the water of reaction according to the invention and deionized water. This is sensible, for example, if owing to the production process there is not sufficient water of reaction available to prepare the desired quantity of dispersion. The aqueous component preferably includes from 0.1 to 80% by weight, and with very particular preference from 1 to 20 by weight, of the water of reaction in the sense of the invention.\nSuitable binders from which it is intended to prepare an aqueous solution or dispersion are in principle all binders known for this purpose. Particularly suitable binders are: polyester, alkyd, acrylate, epoxy, polyurethane and/or epoxide-modified polyester resins, and also modifications of these binders. The preparation of epoxide-modified polyester resins, for example, takes place in accordance with techniques which are known per se.\nThe novel water-dilutable coating compositions may in principle contain all binders suitable for water-dilutable coating compositions.\nBinders which can be used are, accordingly, both refined natural products, for example from rosins or oils or cellulose nitrates, and resins built up entirely by synthetic means The latter include phenolic resins, amine resins (e.g. urea resins, melamine resins), alkyd resins, polyvinyl acetates, epoxy resins, polyurethane resins, polyester resins, rosin-modified phenolic resins, chlorinated rubbers, chlorinated polypropylene, cyclized rubbers, ketone resins and acrylate resins.\nAmino resins, polyester resins, polyacrylate resins and polyurethane resins which can be prepared in organic solution and which are water-dispersible or water-dilutable, and mixtures thereof, are employed in particular as binders.\nThe polyurethane resins employed as binders are known in principle. Examples of suitable polyurethane resins are those described in the literature for use in waterborne coating materials, provided these polyurethane resinsxe2x80x94in a modification of the preparation described in the respective literature can be prepared in the form of organic solutions.\nExamples of suitable polyurethane resins are the resins described in the following documents: EP-A-355433, DE-A 3545618, DE-A 3813866 and the as yet unpublished German Patent Application DE 4005961.8. For further details of the preparation of the polyurethane resins and examples of suitable compounds, therefore, reference may be made to these documents.\nThe polyacrylate resins employed as binders are likewise known and are described, for example, in DE-A 3832826. Suitable polyacrylate resins are General water-dilutable and/or water-dispersible polyacrylate resins which can be prepared in the form of organic solutions.\nAlso suitable as binders are polyester resins which are water-dilutable or water-dispersible and can be prepared in the form of organic solutions. Use is made, for example, of corresponding commercially available water-dilutable or water-dispersible polyester resins, and of the polyester resins which are customarily employed in waterborne coating materials.\nWater-dilutable and water-dispersible amino resins are also suitable as binders. It is preferred to employ water-dilutable melamine resins. Generally, these are etherified melamine-formaldehyde condensation products.\nThe water-solubility of the amino resins dependsxe2x80x94apart from the degree of condensation, which should be as low as possiblexe2x80x94on the etherification component, with only the lowest members of the alcohol or ethylene glycol monoether series giving rise to water-soluble condensation products. The methanol-etherified melamine resins are of greatest importance. If solubilizers are used, butanol-etherified melamine resins can also be dispersed in aqueous phase. Another possibility is to incorporate carboxyl groups into the condensation product. Transetherification products of highly etherified formaldehyde condensation products with hydroxycarboxylic acids are water-soluble via their carboxyl groups, after neutralization, and may be present in the base paints.\nAs binders it is of course also possible to employ mixtures of the binders mentioned and, in addition or alone, other water-dilutable or water-dispersible binders.\nThe novel coating compositions may include customary additions, such as solvents, fillers, plasticizers, stabilizers, wetting agents, dispersion auxiliaries, leveling agents, antifoams and catalysts, and also additives, individually or in a mixture in the customary quantities. These substances may be added to the individual components and/or to the overall mixture.\nExamples of suitable fillers are talc, mica, kaolin, chalk, quartz flour, absestos flour, slate flour, barium sulfate, various silicic acids, silicates, glass fibers, organic fibers or the like.\nIn addition to water, the novel coating compositions may include the customary solvents, for example aliphatic or aromatic hydrocarbons, mono- or polyhydric alcohols, ethers, esters, glycol ethers and their esters, ketones, such as, for example, toluene, xylene, butanol, ethyl- or butylglycol (=ethylene glycol monoethyl or -butyl ether) and their acetates, butyl diglycol (ethylene glycol dibutyl ether), ethylene glycol dimethyl ether, diethylene glycol dimethyl ether, cyclohexanone, methyl ethyl ketone, acetone, isophorone or mixtures thereof.\nIt is additionally possible to add crosslinked polymeric microparticles, as are disclosed, for example, in EP-A-38 127, and/or customary inorganic or organic additives to the novel basecoats. Thus examples of effective thickeners are water-soluble cellulose ethers, such as hydroxyethyl cellulose, methylcellulose or carboxymethylcellulose, and also synthetic polymers having ionic and/or associative groups, such as polyvinyl alcohol, poly(meth)acrylic acid, polyvinylpyrrolidone, styrene-maleic anhydride or ethylene-maleic anhydride copolymers and their derivatives, or else hydrophobically modified ethoxylated piperidine and triethanolamine used [sic]. With particular preference, tertiary amines are employed as neutralizing agents, especially dimethylethanolamine, triethylamine, tripropylamine and tributylamine.\nThe novel basecoat compositions may include all dyes or pigmets which are known and are customary in the coatings industry.\nExamples of dyes and pigments, which may be inorganic or organic in nature, are titanium dioxide, graphite, carbon black, zinc chromate, strontium chromate, barium chromate, lead chromate, lead cyanamide, lead silicochromate, zinc oxide, cadmium sulfide, chromium oxide, zinc sulfide, nickel-titanium yellow, chromium-titanium yellow, iron oxide red, iron oxide black, ultramarine blue, phthalocyanine complexes, naphthol red, quinacridones, halogenated thioindigo pigments or the like.\nAs particularly preferred pigments use is made of metal powders, individually or in a mixture, such as copper, copper alloys, aluminium and steel, preferably aluminium powders, in at least the predominant proportion, namely in a quantity of from 0.5 to 25% by weight based on the overall binder solids content of the coating compositions. Preferred metallic pigments are those commercially available metal powders which have been specially pretreated for aqueous systems. The metal powders may also be employed together with one or more of the abovementioned nonmetallic pigments and/or dyes. In this case their proportion is chosen such that the desired metallic effect is not suppressed.\nThe novel binder solution or binder dispersion is preferably used for producing aqueous coating materials. For this purpose it is processed further, together with crosslinkers (e.g. amino resins, especially those of the type hexamethoxymethylmelamine, urea resin, water-dilutable epoxy resin and the like) optionally further binders, pigments, fillers and other customary additives, to form an aqueous coating material. In addition to water, the latter may also, to a minor extent, comprise organic solvents as diluents.\nThe scope of the present invention also includes the aqueous coating material obtained by using the described binder solution or binder dispersion. The present invention started out, inter alia, from the object of avoiding the problems associated with the water of reaction which occurs in the case of condensation reactions. In the prior art, this water of reaction is separated off from the reaction mixture and then disposed of expensively. The invention, by contrast, proposes employing the water of reaction as an aqueous component in the preparation of aqueous dispersions or solutions of organic substances. Examples of such substances and/or of such a use is the above-described preparation of aqueous binder dispersions and binder solutions. However, it is obvious that the use of the water of reaction is not restricted to these specific groups of substances. Rather, the water of reaction can be employed as solvent whenever the contaminants which it may still contain are harmless for the dispersion or solution to be prepared. The advantage of such a water of reaction is that its composition and the possible accompanying substances which are present in it are known with great precision, provided one knows the precursors of the underlying condensation reaction. In particular, the water of reaction has none of the inorganic impurities of customary tap water.\nThe invention also relates to a process for preparing binder solutions or binder dispersions for aqueous coating materials, which is characterized by the following steps:\na) water of reaction eliminated during the condensation of alcohols and carboxylic acids is separated off from the condensation mixture,\nb) if desired, the water of reaction separated off is mixed with deionized water, and\nc) the binder is dissolved or dispersed in the aqueous component obtained in accordance with a) and b).\nThe water of reaction in process step a) may in principle come from a condensation reaction which has nothing to do with the binder concerned. Preferably, however, at least part of the binder is prepared by the condensation reaction of step a). This is because a particularly rational process is one in which the water of reaction of a polyester or alkyd resin which is employed in aqueous coating systems is employed directly by decantation from the separating receiver in the course of preparing the aqueous colloidal solution of the polyester concerned in a dilution vessel which in most cases is connected to the reactor.\nIn a further, independent process for preparing binder solutions or binder dispersions for aqueous coating materials, the following steps are carried out:\na) binder is prepared by condensation of polyol and poly- or monocarboxylic acids, with the water of reaction eliminated during the condensation remaining in the condensation mixture,\nb) if desired, deionized water is added, and\nc) if desired, further binders are added.\nThis process can be employed whenever at least some of the binder is prepared in a condensation reaction of polyols and carboxylic acids and when the water of reaction which occurs in this condensation reaction, including the impurities which may still be present, does not impair the dispersion which is to be prepared. In this case the process described is particularly advantageous, since it avoids the step of separating off the water of reaction from the condensation mixture. Since the water of reaction in some cases is not present in sufficient quantities for the dispersion to be prepared, it is possible if required to add further deionized water."}
{"text": "1. Technical Field\nThe present invention is generally related to media storage boxes and, more particularly, to lockable media storage boxes.\n2. Background Information\nRenting items of recorded media such as video cassettes and video games has become immensely popular in recent times given the ever increasing number of items available for home viewing and use as well as the decrease in price of the machines that play the media. Typical rental stores display the items available for rental in storage boxes that protect the items from dust, ultraviolet light, and damage from impact if accidentally dropped. Stores protect themselves against theft by placing one or more EAS tags on or inside the item of recorded media. An EAS tag is adapted to activate an alarm when passed through a sensing device that may be disposed around the exit of the store. It is generally desirable to place the EAS tags in places where they cannot be easily removed by a shoplifter.\nA rental establishment typically places older rental stock out for sale when it no longer needs the item for rental. When items are put out for sale, the EAS tags are removed from the items and reused on items that are being rented. The EAS tags are removed from the items of recorded media prior to sale so that the purchaser will not activate the alarm systems in other stores with the EAS tags. Without the EAS tags in the items themselves, a rental store must protect itself against shoplifting by placing an EAS tag on the storage container holding the item for sale. It is desirable that the EAS tag be placed on the storage container in a way that allows a sales clerk to quickly and easily remove the tag after the sale is made but also in a manner that prevents a shoplifter from easily removing the tag. Placement of the EAS tag in such a position is, however, difficult because a storage container provides few, if any, areas where the EAS tag may be hidden. If the EAS tag is placed on the storage container in a fashion such that it can be easily removed, a shoplifter may simply remove the tag and steal the item without activating the alarm. It is thus desired in the art to provide a storage container suitable for holding and displaying an item of recorded media during display for rental that has locking holes that can later accept a lock containing an EAS tag when the container is used to sell an item.\nSuch a storage box and lock combination must be configured to prevent the thief from simply breaking the lock off the storage box with a small pry bar. Although the storage box and lock must be relatively strong, the cost of manufacturing the lock and storage box must not be prohibitive. It is thus also desired in the art to provide a relatively inexpensive storage box and lock and key combination that securely locks the storage box while preventing a pry bar from being inserted between the lock and the storage box."}
{"text": "The present invention relates to an electrostatic protection circuit and a semiconductor integrated circuit using the same, particularly to an electrostatic protection circuit for an integrated circuit using an insulated gate field effect transistor (hereinafter referred to as a MOSFET).\nFIG. 9 shows a typical electrostatic protection circuit having a commonly-used signal terminal 1509 and power supply terminals 1 and 2 of a MOS integrated circuit using a bulk substrate. In FIG. 9, there are two paths for absorbing static electricity charge applied between the signal terminal 1509 and the power supply terminals 1 and 2. One is a path in which discharge current flows from the signal terminal 1509 to the first power supply terminal 1 designated as potential +VDD through a diode 1503 and the other is a path in which discharge current flows from the second power supply terminal 2 designated as potential xe2x88x92VSS to the signal terminal through a diode 1504. There is another path in which discharge current flows from the second power supply terminal 2 through the diode 1501 to the first power supply terminal 1. To be more practical, an input signal from a pad terminal 1506 is fed to gate electrodes of a p-type MOSFET 1507 and an n-type MOSFET 1508 which form an internal circuit inverter, through a resistor 1505 and the connection of the terminals of the diodes 1503 and 1504, as shown in FIG. 9.\nIf the pad 1506 is directly connected to the gate of the p-type MOSFET 1507 or n-type MOSFET 1508 of the internal circuit in FIG. 9, the gate electrodes of the p-type MOSFET 1507 and the n-type MOSFET 1508 are often broken when static electricity is applied from the pad 1506. To prevent this, a resistor 1505 for buffering the shock of static electricity and the diodes 1503 and 1504 for absorbing the charge are used. In addition, the diode 1501 serves as a charge-absorbing path not only for static electricity applied between the first and second power supply terminals but also for static electricity applied to the above-described signal terminals as discussed later.\nIn the conventional electrostatic protection circuit, the above-mentioned diode element 1503 for absorbing charges is connected so as to conduct the charges to the first power supply terminal 1, and the diode element 1504 is connected so as to conduct the charges from the second power supply terminal 2 to the signal terminal 1509. Further, the diode 1501 is connected in a reverse direction between the first power supply terminal 1 and the second power supply terminal 2. This is because, if the diodes 1501, 1503, and 1504 are connected in a forward direction so as to conduct the current in the opposite direction to that shown in FIG. 9, leakage currents may flow through the forward-biased diodes when the integrated circuit is connected to a power supply.\nFurther, in integrated circuits using a silicon-on-insulator substrate (hereinafter abbreviated as xe2x80x9cSOI integrated circuitxe2x80x9d), there are no wells such as those in a bulk substrate, the bottom is insulated by a buried oxide film, and the side is also covered with a local oxide film formed by LOCOS (local oxidation of silicon) method. For this reason, there is generally no equivalent to the diode 1501 between the first and second power supply terminals shown in FIG. 9. In other words, there is no diode formed by a p well and an n well in a conventional substrate between the first and second power supply terminals as shown in FIG. 11. However, such a diode is required from the viewpoint of electrostatic protection as discussed later. Therefore, in an SOI integrated circuit, a diode 1801 is added between the first and second power supply terminals as shown in FIG. 12, or a MOSFET 1901 with the source and gate electrodes connected together is connected between the first and second power supply terminals for causing the MOSFET 1901 to perform the function of the reverse-biased diode as shown in FIG. 13. Alternatively, a p-type MOSFET 2001 and n-type MOSFET 2000 each having the source and gate electrodes connected together are connected in parallel between the first and second power supply terminals as shown in FIG. 14 to cause them perform the function of the reverse-biased diodes. SOI integrated circuits are thus provided with electrostatic protection based on the same principle as for bulk-substrate integrated circuits.\nIf static electricity is added between the power supply terminals or between the signal terminal and one of the power supply terminals, electrostatic breakdown may occur inside the integrated circuit. In the case where static electricity is added to signal terminals, the internal circuit to which the signal terminals are connected may frequently be damaged unless the charges are quickly absorbed by an electrostatic protection circuit. In FIG. 9, the gates of p-type MOSFET 1507 and n-type MOSFET 1508 may be broken down. The gate insulating film of a MOSFET is made of a very thin film from several hundred angstrom to several tens angstrom formed between the substrate and the gate electrode, and the substrate or the source electrode is finally connected to the power supply. As a result, a high voltage is applied across the thin gate silicon oxide film to produce a strong electric field, resulting in breakdown of the gate film. Therefore, when static electricity is added, the electrostatic protection circuit in FIG. 9 or a similar means is used to absorb the charges quickly. In the circuit of FIG. 9, the following four cases are possible for the flow of charges between the signal terminal 1509 and the first and second power supply terminals 1 and 2:\n(A) Signal terminal: positive charges, First power supply terminal: negative charges\n(B) Signal terminal: negative charges, First power supply terminal: positive charges\n(C) Signal terminal: positive charges, Second power supply terminal: negative charges\n(D) Signal terminal: negative charges, Second power supply terminal: positive charges\nIn the conventional circuit shown in FIG. 9, the diode 1503 or diode 1504 operates in a forward direction in the cases (A) and (D) above. Accordingly, electrostatic charges that comes in can be quickly absorbed, thereby preventing electrostatic destruction. In cases (B) and (C), the diodes 1503 and 1504 are both in a reverse direction with respect to the polarity of electrostatic charges. In case (B), negative charge forces its path through the diode 1503 in a reverse direction. Otherwise, the negative charges first flow through the diode 1504 in a forward direction, then flow from the second power supply terminal 2 to the first power supply terminal 1 through the reverse diode 1501 in the substrate. In case (C), positive charges force its path through the diode 1504 in a reverse direction. Alternatively, the positive charges first flow through the diode 1503 in a forward direction, then flow from the first power supply terminal 1 to the second power supply terminal 2 through the reverse diode 1501 in the substrate. Therefore, in cases (B) and (C), because charges must necessarily flow through a diode in a reverse direction, the circuit is vulnerable to static electricity and can be broken by even comparatively low voltages. A path through which electrostatic charges flow out in the case (C) is shown in FIG. 10 as an example.\nThe function of the reverse-biased diode between the power supplies in the case where static electricity is applied to the signal terminals has been described above. If static electricity is applied between the power supply terminals in the same polarity as the power supply voltage, the charges flow through the diode between the power supply terminals in a reverse direction. Since even this diode does not exist in SOI integrated circuits, electrostatic charges cannot be absorbed, and hence breakdown easily occurs at the most vulnerable part among the parts involved in the power supply lines.\nNext, the mechanism why a diode is endurable in a forward direction and weak in a reverse direction will be briefly discussed referring to FIGS. 30 and 31. In FIG. 30, a p-type diffusion layer 131 and an n-type diffusion layer 132 join together at the boundary to form a pn diode. FIG. 30 shows the case in which the p-type diffusion layer 131 is at a positive potential and the n-type diffusion layer 132 is at a negative potential. In this case, since the diode is forward-biased and a current therefore can flow easily everywhere in the pn boundary, the current flows uniformly across the entire pn boundary. Accordingly, the diode as a whole allows easy flow of current and exhibits high charge-absorbing capability. Furthermore, because the current flows in the diode itself well dispersed and uniformly without extreme concentration, no destruction of the diode itself occurs due to flow of current.\nOn the other hand, FIG. 31 shows the state where a current flows in a reverse direction. In FIG. 31, a p-type diffusion layer 141 and an n-type diffusion layer 142 join together at the boundary to form a pn diode. FIG. 31 shows the case in which the p-type diffusion layer 141 is at a negative potential and the n-type diffusion layer 142 is at a positive potential.\nIn this case, since the diode is reverse-biased to the voltage applied, normally no current flows. However, if a high voltage is applied to the diode in order to force a current in a reverse direction, the current begins to flow at specific parts where the breakdown voltage is comparatively low and hence current can easily flow due to non-uniform pn boundaries. Therefore, even if the voltage goes over the breakdown voltage and a current begins to flow, this does not mean that the current is flowing uniformly across the entire boundary. The current tends to concentrate in the parts allowing easy flow thereof. The state of this current flow is shown in FIG. 31. Since current does not flow uniformly in the diode but tends to concentrate in specific parts when the current flows in a reverse direction in the diode, the capability of the diode to allow current to flow is small relative to the size of the pn boundary area of the diode. In other words, the diode exhibits only a small charge absorbing capability. In addition, since current concentrates in parts of the diode allowing easy flow of current, the current density in those parts extremely increases, there is a high risk for the diode to be unduly heated and broken down.\nAs described above, a diode exhibits high endurance to static electricity when acting in a forward direction, but exhibits so weak endurance when acting in a reverse direction that the diode may be damaged even by low static electricity depending on the manner in which the static electricity is applied.\nIn the conventional electrostatic protection circuit configuration, the above-described occurrence in which the diodes operate only in a reverse direction to the charges is unavoidable. To improve the electrostatic endurance in these cases, conventionally the size of the diode or MOSFET has been increased to cope with weakness in a reverse direction by increasing the areas of the diode or MOSFET. This increases the area occupied by the electrostatic protection circuit attached to each pad, resulting in a cost increase and limitation to the number of pads and pins which can be formed.\nFurthermore, the increase in the diode area accompanies an increase in the parasitic electrostatic capacity acting as a capacitor, causing degradation of the high frequency characteristics of signal terminals for which the high frequency operation is required and an increase in power consumption.\nIn the case of SOI integrated circuits, as mentioned above there is a problem of very low electrostatic endurance due to absence of even the reverse-biased diode formed by p and n wells between the power supply terminals. Assume that an attempt is made, for example, to add a reverse-biased diode or equivalent device corresponding to the reverse-biased diode formed by p and n wells in a conventional bulk substrates between the power supply terminals by forming an additional diode or MOSFET. This method also cannot provide SOI integrated circuits with a sufficiently high electrostatic endurance because it is not possible to form a reverse-biased diode as large as the conventional reverse-biased diode formed by the p and n wells parasiting between power supply terminals with a very large area because of the limitation to the area.\nThe present invention has been completed in view of these problems, and has an objective of providing an improved electrostatic protection circuit which affords a sufficiently high electrostatic endurance between the power supply terminals or between one of the power supply terminals and the signal terminal, and a semiconductor integrated circuit using the electrostatic protection circuit. Another objective of the present invention is to provide an electrostatic protection circuit exhibiting high electrostatic endurance with a comparatively small area, and a semiconductor integrated circuit using this electrostatic protection circuit, which is sufficiently endurable to the static electricity and suitable for high-frequency operation because of very small parasitic electrostatic capacity due to the use of this protection circuit.\nAccording to one aspect of the present invention, there is provided an electrostatic protection circuit comprising:\na first power supply terminal to which a first voltage is applied;\na second power supply terminal to which a second voltage lower than the first voltage is applied;\na first diode connected in a reverse direction between the first and second power supply terminals; and\na second diode connected in a forward direction between the first and second power supply terminals,\nwherein a forward drop voltage of the second diode is set to be higher than a drive voltage supplied between the first and second power supply terminals.\nAccording to this aspect of the present invention, either one of the first and second diodes always operates in a forward direction with respect to the static electricity applied between the first and second power supply terminals regardless of the polarity of the static electricity. Electrostatic charges therefore can be quickly absorbed through the diode in a forward direction.\nBecause the second diode is connected in a forward direction between the first and second power supply terminals in this embodiment, this second diode is forward-biased to the electric charges of the drive voltage applied between the first and second power supply terminals in the ordinary operations. However, the forward drop voltage of the second diode is set to be higher than the drive voltage supplied between the first and second power supply terminals. For this reason, a forward leakage current does not flow through the second diode during normal operation.\nThe second diode may have a pn junction structure formed by a p-type diffusion layer and an n-type diffusion layer joined together. In this case, the forward drop voltage of the second diode is defined by a contact potential (or contact potential difference) which is a potential difference generated at the boundary of the p-type diffusion layer and an n-type diffusion layer. Therefore, if the forward drop voltage of the second voltage is set to be higher than the drive voltage supplied between the first and second power supply terminals, a forward leakage current of the second diode can be prevented from flowing during normal operation.\nThe second diode may be formed from a plurality of diodes connected in series. Each of the plurality of diodes has a pn junction structure formed by a p-type diffusion layer and an n-type diffusion layer joined together. Assuming that the number of diodes connected in series is n, the forward drop voltage of the second diode becomes n times the contact potential of each diode. As a result, it is possible to use higher power supply voltages.\nThe second diode may include first p-type and n-type diffusion layers and a second p-type or n-type diffusion layer disposed between and connected to the first p-type and n-type diffusion layers. In this case, it is preferable that the diffusion concentration of the first p-type and n-type diffusion layers is set to be higher than the diffusion concentration of the second p-type or n-type diffusion layer. Because the contact potential of the second diode can be increased in this manner, it is possible to use higher power supply voltages.\nThe second diode may also be formed by a MOS transistor having a drain electrode and a gate electrode which are connected to each other. In this case, the forward drop voltage of the second diode is defined by the threshold voltage of the MOS transistor.\nThe second diode may also be formed by connecting a plurality of MOS transistors in series. In this case, a drain electrode and a gate electrode are connected to each other in each of the plurality of MOS transistors. Assuming the number of MOS transistors connected in series is n, the forward drop voltage of the second diode is n times the threshold voltage of each MOS transistor. As a result, it becomes possible to use higher power supply voltages.\nThe first diode may be formed by connecting in parallel a p-type MOS transistor having a source electrode and a gate electrode connected to the first power supply terminal, with an n-type MOS transistor having a source electrode and a gate electrode connected to the second power supply terminal. This configuration ensures the first diode to exhibit more stable characteristics.\nThe first and second diodes may be formed on a silicon-on-insulator (SOI) substrate. The first and second diodes are then surrounded by insulator layers, and therefore formation of superfluous parasitic diodes can be prevented.\nAccording to another aspect of the present invention, there is provided an electrostatic protection circuit comprising:\na first power supply terminal to which a first voltage is applied;\na second power supply terminal to which a second voltage lower than the first voltage is applied;\na signal terminal to which a signal voltage equal to or lower than the first voltage and equal to or higher than the second voltage is applied;\na first diode connected in a forward direction between the first power supply terminal and the signal terminal;\na second diode connected in a forward direction between the signal terminal and the second power supply terminal;\na third diode connected in a reverse direction between the first power supply terminal and the signal terminal; and\na fourth diode connected in a reverse direction between the signal terminal and the second power supply terminal,\nwherein a forward drop voltage of each of the first and second diodes is set to be higher than a drive voltage supplied between the first and second power supply terminals.\nAccording to this aspect of the present invention, one of the first to fourth diodes certainly operates in a forward direction with respect to the static electricity applied between the signal terminals and the first/second power supply terminal regardless of the polarity of the static electricity. Electrostatic charges therefore can be quickly absorbed through the diode in a forward direction.\nThe first and second diodes are connected in a forward direction between the signal terminal and the first or second power supply terminal. Therefore, the first and second diodes are forward-biased to electric charges of the drive voltage applied between the signal terminal and the first/second power supply terminal in the ordinary operations. However, the forward drop voltage of each of the first and second diodes is set to be higher than the drive voltage supplied between the first and second power supply terminals. For this reason, a forward leakage current does not flow through the first and diode during normal operation.\nThe above-mentioned features can be applied to other aspects of the present invention.\nAccording to still another aspect of the present invention, there is provided a semiconductor integrated circuit comprising:\na logic circuit formed by interconnecting a plurality of p-type MOS transistors and a plurality of n-type MOS transistors; and\nan input/output circuit disposed around the logic circuit,\nwherein the input/output circuit includes an electrostatic protection circuit which protects the logic circuit from static electricity and has elements of one aspect of the present invention.\nAccording to this aspect of the present invention, the logic circuit can be protected from static electricity by the actions of the above-described electrostatic protection circuit.\nAbove-described features of the electrostatic protection circuit can be applied to the semiconductor integrated circuit of the present invention.\nParticularly, the diffusion concentration of at least one of the p-type and n-type diffusion layers of the first diode may be set equal to or higher than the diffusion concentration of diffusion layers used for source electrodes of the plurality of p-type and n-type MOS transistors in the logic circuit. If the diffusion concentrations of the above diffusion layers are set to be equal, no forward leakage current occurs through the second diode during normal operation when the forward drop voltage (contact potential) of the second diode is set to be higher than the drive voltage supplied between the first and second power supply terminals. If the above diffusion concentrations are set to be higher, the contact potential of the second diode increases and higher power supply voltages can be used.\nSimilarly in the case that the threshold voltage of the MOS transistor which forms the second diode is set to be higher than the threshold voltage of each of the plurality of p-type and n-type MOS transistors in the logic circuit, higher power supply voltages can be used.\nThe semiconductor integrated circuit according to other aspects of the present invention includes the electrostatic protection circuit having the elements of the other aspects of the present invention.\nThe above described features of the electrostatic protection circuit can be applied to the semiconductor integrated circuit of the present invention."}
{"text": "Japanese Patent Application Laid-Open Publication No. 2008-118004 (Patent Document 1) describes a layout structure in which a power supply wiring connected to a diffusion layer is drawn onto boundary line between adjacent standard cells. More specifically, Patent Document 1 describes a layout structure in which a diffusion layer is shared between adjacent standard cells and a wiring is drawn from a power supply wiring onto boundary line between the adjacent standard cells. And, Patent Document 1 describes a structure in which the wiring drawn on the boundary line between the adjacent standard cells and the diffusion layer shared between the adjacent standard cells are electrically connected to each other via a plug."}
{"text": "Steam bath have been popular worldwide recently for health reasons. The only drawback of a steam bath is the high service charge due to the expensive associated facilities. As such, a relatively cheap portable sauna enclosure has been commercially available as shown in FIG. 1. The sauna enclosure comprises a frame 1 including a plurality of interconnected vertical and horizontal bars, a canopy 2 covered on the frame 1 and stretched by it for forming a space to allow a person to take a bath therein, canopy 2 having three holes 201 for the neck and two arms extending through and a zipper 202 for permitting a person to enter the sauna enclosure when it is unfastened, a steamer 3 outside the sauna enclosure, and a steam pipe 4 connected between the steamer 3 and the inside of the sauna enclosure for transferring steam from steamer 3 to the sauna enclosure. The user should add water into steamer 3 for heating. After bathing, the user should use a dry cloth to absorb the condensed steam (i.e., water) in the sauna enclosure in order to forestall the growth of micro-organisms on the canopy 2.\nHowever, the previous design suffered from several disadvantages. For example, it is still possible for steam to pass through holes 201 and zipper 202 during bathing due to the gaps between human body (i.e., neck and arms) and holes 201 as well as the gaps in the zipper 202 itself. As such, the leaked steam may be subsequently carried from the bathroom to the other rooms in the house, resulting in a potential damage to the household appliances, furniture, and wall due to humidity. In addition, the assembly and disassembly of frame 1 and steam pipe 4 is a tedious job. Also, the efficiency of steamer 3 may be adversely affected if the packing between steam pipe 4 and steamer 3 is not well sealed. Furthermore, the bulky steamer 4 is not easy to store when not in use. In view of this, its portability is adversely affected. Thus, it is desirable to provide an improved portable sauna enclosure in order to overcome the above drawbacks of the prior art."}
{"text": "Brine is a concentrated salt solution (usually sodium chloride), which is used for preserving various foodstuffs, as well as in the refrigerant industry. Brine is also a by-product of some seawater desalination plants. However, the disposal of brine can be problematic. For example, the density of brine is greater than typical seawater, and if it was discharged into the sea, it would tend to sink towards the seabed and not disperse throughout the water column. This can result in a hypersaline layer of water forming on the seabed, which could damage sea habitats and organisms."}
{"text": "This invention relates to a cyclic or step feed drilling machine in which one drilling process is carried out, instead of being performed in a single stroke, by gradually deepening a bore-depth by means of repeating reciprocative movement of a tool for a plurality of times before completing the drilling process.\nMore particularly it relates to a cyclic drilling machine in which magnitude of a load applied on the tool is detected so that the tool is, when the detected load grows larger to a certain value, retreated backwards once for being advanced again.\nAs a prior art belonging to this category, a technological idea disclosed in the Japanese patent publication TOKU-K0-No. SHO 53(1978)-40745 can be cited, in which magnitude of a load applied on the tool is detected or measured by the fluctuation or variation of load current observed in a tool driving motor, and the tool, i.e. drill in this instance, is returned backwards once in response to an indication of reaching of the load current to a predetermined value.\nIn this kind of device a load current of a motor at a non-load status is memorized once as a minimum value (Io) in a memory circuit, and a load current of a motor corresponding to a cutting (drilling) torque within the maximum allowable limit, wherein the drill can not be broken, is set in advance on the other hand as a current set value ({ Is). Sum of the both values (Is (=Io+.DELTA.Is)) is regarded as a reference value. Comparing of this reference value (Is) with a load current (I) of a motor which is detected in the course of a cutting is executed in a comparing circuit or comparator, so that the drill in advancing movement may be instantly returned backwards as soon as the load current (I) reaches the reference value (Is), that is when a cutting torque of the maximum allowable limit is applied on the drill. When a certain predetermined period of time has elapsed after the motor was placed under a non-load condition due to the return of the drill, the then non-load current value is memorized in the memory circuit in place of the previous minimum value as a new minimum value (Io). The drill is advanced again so as to resume the cutting of a workpiece which has been halfway suspended. The then load current value (I) of the motor is compared in the comparator with the reference value (Is). When the drill receives a cutting torque of the maximum allowable limit again it is instantly returned backwards similarly to the previous instance. This type of operation is repeated until a bore of a predetermined depth is formed in the workpiece, so that the drill may be prevented from damage or breakage.\nThe cutting process is however finished in general within 1-3 minutes after a workpiece is placed on the drilling machine, and the motor is stopped between the finish of cutting on one workpiece and the beginning of cutting on another workpiece. In other words, the cutting is carried out in the prior art while the motor (in this instance AC motor is used, but DC motor is also permissible) is under an unstable status, as shown in FIG. 1, wherein the winding current gradually decreases from a large initial current value to a stable current value. It means that the aforesaid minimum value (Io) is set while the winding current is unstable in the course of falling as shown in FIG. 2 and the reference value (Is), on which the cutting process depends, is set according to this minimum value (Io). Comparison of the minimum value (Io) with a current value (Io1), which would have been attained if the minimum value had been set when the winding current reached the reference value (Is), will show that the current value (Io1) is smaller than the minimum value (Io) by an amount of (.DELTA.I.alpha.).\nFor preventing the drill from the breakage, the reference value (Is) must be made, therefore, smaller by the difference (.DELTA.I.alpha.) of the minimum value (Io) minus the current value (Io1), i.e., to (Iss (=Is-.DELTA.I.alpha.)). It signifies that for making the reference value (Is) smaller by .DELTA.I.alpha., the aforementioned current set value (.DELTA.Is) must be set smaller by that much. In practical machining lines, such as a machining factory, wherein power source is variable in voltage even when a motor itself is in a stable status, the winding current thereof varies a great detal so that mistaking of this variation for that owing to load variation will cause an error operation. External noise will also cause the error operation. Reducing of the current set value (.DELTA.Is) is therefore envisaged with a limit. The step feed type cutting (drilling) with a thin drill which requires setting of a small current set value (.DELTA.Is) is practically next to impossible.\nComparison of a characteristic curve Ip1 of the winding current in the initial stage of the motor driving starting from a sleeping status with those curves Ip2 and Ip3 in the initial stage of the motor driving at a repeatedly operated status will show that the time required for decreasing value is longer in the curve Ip1 than in the others Ip2, Ip3, due to variation in the motor itself such as internal resistance in the motor. And the difference (I.alpha.) is maintained in a larger state, in other words, a status where the difference (.DELTA.I.alpha.) is large is maintained for a long period of time in the drilling of the workpiece in the initial operation period. It means that the time with a large difference (.DELTA.I.alpha.) continues for a comparatively longer period, particularly in a colder region for quite a long time, before the motor enters a stabilized operation wherein the machining can be stably performed. It naturally deteriorates the operational efficiency, to the great disadvantage.\nIn the conventional devices variation of the load current, i.e. variation of the winding current of the motor which is caused by the load variation is the detection object. It is quite difficult to get true variation of the cutting torque in this method, because the variation obtained herewith is nothing but a sum of the true variation value of the cutting torque plus a torque produced in the transmission system from the motor to the drill. And the torque produced in the transmission system is variable in itself during operation thereof. It is, therefore, problematical in this connection whether the variation of the winding current can be truly regarded as the variation of the cutting torque.\nBesides, the technology disclosed in the Japanese patent publication TOKU-SO-No. SHO 54(1979)-23158 attempts to raise in put impedance by means of an emitter follower with FET (field-effect-transistor) so as to minimize leakage of electric charge which is charged in a capacitor for the purpose of preventing the timewise variation of the minimum value (Io) in a memory circuit storing the minimum value (Io). It is however practically very difficult to prevent the leakage in the capacitor, so exact maintenance of the minimum value (Io) and subsequent setting as well as maintaining of an exact reference value (Io) is next to impossible, to the great problem."}
{"text": "Nipple drinker watering systems are commonly used to provide water to poultry and other small animals. These watering systems involve the use of several branching water supply lines extending the interior length of a house, such as a poultry house. The water supply lines have numerous nipple drinkers attached thereto so that the poultry may obtain water by pecking at the nipples. Stand-tubes are typically provided at various locations along the length of the water supply lines, typically proximate to pressure regulators. The stand-tubes are generally clear and provide a visual indication, with the assistance of a float provided in the stand-tube, of the water pressure level (via water column) in the water supply line at their locations. In order to be effective, the stand-tubes must be able to vent to the outside environment.\nPeriodically, it may be necessary to flush the nipple drinker watering systems in order to clean out the watering system as material may build up in the watering lines over time due to, for instance, the introduction of sediment from the water source itself or from growers adding material to the water source, such as chlorine or medication. Such build-up of material in the watering system can interfere with proper water flow along the water supply line and may eventually contribute to the malfunction of the nipple drinkers. It may also be necessary to flush the nipple drinker watering systems in order to remove warmer water from the watering system. Removal of warm water from the watering system in order to replace it with cooler water has been found to increase consumption of water by poultry.\nDue to the need to flush the systems, the stand-tubes have been outfitted with caps designed to try and prevent leakage out of the stand-tubes during a flushing operation, as well as to prevent the floats from being blown out of the stand-tubes in view of the increased water pressure in the system caused by the flushing. One example of such a cap is shown and described in U.S. Pat. No. 5,136,983. This cap, as well as others, however, have had their own problems. For instance, some caps have had problems with air lock being caused by the seal being formed when it should be allowing air to pass. Other caps have had problems with still allowing leakage to occur after the seal is formed. Other caps have had problems with forming the seal at lower water pressures (e.g., the seal may be formed in caps provided closer to the water source, where the water pressure is typically higher, but not in caps provided distal from the water source, where the water pressure is typically lower). Other caps have designs which allow for air flow both in and out of the assembly, but which allows for unwanted particles to enter the assembly which can cause contamination of the water and which can creates blockages in the cap, thereby potentially causing it to malfunction.\nIn view of the foregoing, there is a need for a breather cap assembly which overcomes all of the disadvantages of the prior art caps."}
{"text": "The present invention relates to a stage apparatus used for, e.g., a fine moving stage for finely aligning a substrate such as a wafer or a master such as a reticle in an exposure apparatus or the like, an exposure apparatus, and a device manufacturing method.\nControl techniques for industrial and information devices are increasing in precision and speed as such devices become complicated, advanced, and miniaturized. Exposure apparatuses for manufacturing xe2x80x9csemiconductor devicesxe2x80x9d which support various types of such devices are also rapidly increasing in speed and precision, and demands are arising for nm-order alignment precision and the like.\nA generally known alignment mechanism for a wafer and reticle in an exposure apparatus is a stage apparatus constituted by a fine moving stage having a microtranslation mechanism and fine rotation mechanism, and a coarse moving stage having a large-stroke moving mechanism in the translational direction.\nSuch an exposure apparatus finely aligns, controls, and exposes a wafer or reticle on the fine moving stage so as to make it coincide with the image or object plane by the fine moving stage mounted on the coarse moving stage while the coarse moving stage makes large-stroke movement.\nIn general, both the coarse and fine moving stages often use linear motors, which utilize a Lorentz force, as the actuator of the alignment mechanism.\nHowever, the fine moving stage of the exposure apparatus has recently been found to suffer from the following problems.\n(1) As chips, wafers, and reticles increase in size, the fine moving stage also increases in size. Pushing up the fine moving stage by a linear motor using a Lorentz force consumes a large amount power.\nEspecially the actuator of the fine moving stage is often attached near a wafer or reticle. Compared to the coarse moving stage, the fine moving stage greatly influences the wafer or reticle to thermally deform it, resulting in very low final exposure precision.\n(2) Demands are arising for a high-speed exposure apparatus, and the acceleration of the coarse moving stage is being increased. This means that the acceleration applied from the coarse moving stage to the fine moving stage also increases. When the coarse moving stage accelerates the fine moving stage, a moment force proportional to the acceleration of the coarse moving stage acts on the fine moving stage due to the difference between the barycenter of the fine moving stage and a xe2x80x9cpoint of applicationxe2x80x9d where the actuator of the coarse moving stage pushes the fine moving stage.\nHence, as the speed increases, the actuator of the fine moving stage becomes bulky in order to withstand the moment force, and the amount of heat generated by the actuator increases to adversely influence the exposure precision. In other words, a bulky actuator increases heat, which is a vicious cycle.\nAs a means for solving these problems, there is provided a method of compensating for the weight of the fine moving stage by using the repulsion and attraction forces of a magnet or electromagnet. In this case, it is difficult to hold the posture of the fine moving stage by compensation with only the magnet, so a spring is also used to hold the posture.\nFIGS. 5A and 5B schematically show a conventional wafer stage for compensating for the weight by magnetic repulsion. This wafer stage comprises a fine moving stage 101 having a support surface 101a for supporting a wafer WO, a coarse moving stage 102 having a large-stroke moving mechanism (not shown), first to third spring members 103a to 103c for holding the posture of the fine moving stage 101 on the coarse moving stage 102, and first to third magnetic units 104a to 104c for compensating for the weight by a so-called magnetic repulsion force for supporting the weight of the fine moving stage 101. The magnetic units 104a to 104c are respectively made up of magnets 141a to 141c fixed to the upper surface of the coarse moving stage 102 and magnets 142a to 142c fixed to the lower surface of the fine moving stage 101 so as to face the corresponding magnets 141a to 141c. \nThe weight of the fine moving stage 101 is compensated for by the magnetic repulsion forces of the magnets 141a to 141c and 142a to 142c of the magnetic units 104a to 104c. When the wafer W0 is aligned near the image plane of a projection optical system (not shown), the weight of the fine moving stage 101 is canceled by the magnetic repulsion force.\nThe three spring members 103a to 103c for compensating for the posture of the fine moving stage 101 have the same spring constant, which is set to minimize the restoring forces of the springs while the wafer W0 is near the image plane.\nThe fine moving stage 101 has a pair of bar mirrors 110 for detecting the position of the fine moving stage 101 by a laser interferometer or the like. This wafer stage is equipped with a control system for feeding back the detected position information of the fine moving stage 101 to the actuator of the coarse moving stage 102.\nThe three posture compensation spring members 103a to 103c are generally arranged at positions Z1 to Z3 at equal intervals of 120xc2x0 on a circle of a radius r having a geometrical center O of the fine moving stage 101 as its center.\nIn this prior art, however, the barycenter of a structure formed from the members and fine moving stage is at a position different from the geometrical center of the fine moving stage, i.e., the geometrical centroid of a triangle formed by the support points Z1 to Z3 of the spring members 103a to 103c under the influence of the members such as the bar mirrors fixed to the fine moving stage for holding a wafer. The restoring forces of the three spring members necessary for pushing up the fine moving stage only in the +Z direction without rotating it and for keeping it at a predetermined pushed-up position and posture are different.\nAccordingly, holding the posture by the three springs having the same spring constant causes the following problems even if the magnets compensate for the weight of the fine moving stage. Assume that the fine moving stage is vertically displaced in the xe2x88x92Z direction without rotating it, and the springs uniformly deflect. In this case, the sum of the restoring forces of the three springs not only translates the fine moving stage in the +Z direction but also rotates it.\nMore specifically, if an external force acts on the fine moving stage owing to some factor while position control by the actuator of the fine moving stage stops, or if position control stops due to an emergency stop caused by an accident while the fine moving stage is displaced in the xe2x88x92Z direction by the actuator, restoring motion of the fine moving stage in the +Z direction and tilt motion exceeding a moving stroke in the rotational direction occur, scoring or damaging the stage.\nThere is proposed a method of compensating for the weight of the fine moving stage by magnets, compensating for the posture by springs, and generating a uniform weight compensation force by the magnets in almost the entire range of the Z-direction stroke on the fine moving stage. Even this method requires a posture holding means using coil springs or leaf springs because compensation by only magnets is unstable. This method cannot solve the static unbalance.\nThere is also proposed a method of directly compensating for the weight of the fine moving stage by springs. Even this method cannot solve the above problems.\nThe present invention has been made to overcome the conventional drawbacks, and has as its object to provide a stage apparatus which uses at least three spring members in order to compensate for the weight of a stage for setting a wafer or reticle, can improve the stability of a compensation function by the spring members, and can effectively avoid scoring of or damage to the stage, to provide an exposure apparatus including the stage, and to provide a method of manufacturing a device by using the exposure apparatus.\nTo achieve the above object, according to the present invention, there is provided a stage apparatus comprising:\na stage;\na member fixed to said stage; and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including said stage and said member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined based on a relative position relationship between the barycentric position of the structure and the at least three spring members so as to set at least one spring constant different from the other spring constants.\nAccording to a preferred embodiment of the present invention, it is preferable that the at least three spring members include coil springs.\nAccording to a preferred embodiment of the present invention, it is preferable that the at least three spring members include air cylinders.\nAccording to a preferred embodiment of the present invention, it is preferable that the at least three spring members include leaf springs.\nAccording to a preferred embodiment of the present invention, it is preferable that the weight compensation mechanism further includes a magnet.\nAccording to a preferred embodiment of the present invention, it is preferable that the apparatus further comprises an actuator for vertically driving said stage.\nAccording to a preferred embodiment of the present invention, it is preferable that the actuator includes a linear motor.\nAccording to a preferred embodiment of the present invention, it is preferable that the apparatus further comprises at least three actuators for vertically driving said stage, wherein a supporting force generated by each of the spring members and applied to said stage and a driving force generated by each of the actuators and applied to said stage are exerted on the same axis.\nAccording to a preferred embodiment of the present invention, it is preferable that the member includes a mirror.\nAccording to a preferred embodiment of the present invention, it is preferable that the at least three spring members are arranged at positions a constant distance from the geometrical center of said stage.\nAccording to a preferred embodiment of the present invention, it is preferable that the at least three spring members include n spring members, and the n spring members are arranged at vertexes of a regular n-side polygon having the geometrical center of said stage as a center.\nAccording to the present invention, there is provided a stage apparatus comprising:\na stage;\na member fixed to said stage; and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including said stage and said member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined to make application points of restoring forces of the at least three spring members coincide with the barycentric position of the structure when said stage is statically balanced.\nAccording to the present invention, there is provided a stage apparatus comprising:\na coarse moving stage;\na fine moving stage mounted on said coarse moving stage;\na member fixed to said fine moving stage; and\na weight compensation mechanism which has at least three spring members for compensating for a weight of a structure including said fine moving stage and said member and is interposed between said coarse moving stage and said fine moving stage, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined based on a relative positional relationship between the barycentric position of the structure and the at least three spring members so as to set at least one spring constant different from the other spring constants.\nAccording to the present invention, there is provided a stage apparatus comprising:\na coarse moving stage;\na fine moving stage mounted on said coarse moving stage;\na member fixed to said fine moving stage; and\na weight compensation mechanism which has at least three spring members for compensating for a weight of a structure including said fine moving stage and said member and is interposed between said coarse moving stage and said fine moving stage, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined to make application points of restoring forces of the at least three spring members coincide with the barycentric position of the structure when said stage is statically balanced.\nAccording to the present invention, there is provided an exposure apparatus comprising:\nan illumination system for illuminating a master;\na master stage apparatus which moves while supporting the master;\na projection optical system for projecting a pattern of the master onto a substrate; and\na substrate stage apparatus which moves while supporting the substrate,\nat least one of said master and substrate stage apparatuses having\na stage,\na member fixed to said stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including said stage and said member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined based on a relative positional relationship between the barycentric position of the structure and the at least three spring members so as to set at least one spring constant different from the other spring constants.\nAccording to the present invention, there is provided an exposure apparatus comprising:\nan illumination system for illuminating a master;\na master stage apparatus which moves while supporting the master;\na projection optical system for projecting a pattern of the master onto a substrate; and\na substrate stage apparatus which moves while supporting the substrate,\nat least one of said master and substrate stage apparatuses having\na stage,\na member fixed to said stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including said stage and said member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined to make application points of restoring forces of the at least three spring members coincide with the barycentric position of the structure when said stage is statically balanced.\nAccording to the present invention, there is provided an exposure apparatus comprising:\na projection optical system for projecting an image onto a substrate; and\na substrate stage apparatus which moves while supporting the substrate,\nsaid substrate stage apparatus having\na stage,\na member fixed to said stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including said stage and said member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined based on a relative positional relationship between the barycentric position of the structure and the at least three spring members so as to set at least one spring constant different from the other spring constants.\nAccording to the present invention, there is provided an exposure apparatus comprising:\na projection optical system for projecting an image onto a substrate; and\na substrate stage apparatus which moves while supporting the substrate,\nsaid substrate stage apparatus having\na stage,\na member fixed to said stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including said stage and said member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined to make application points of restoring forces of the at least three spring members coincide with the barycentric position of the structure when said stage is statically balanced.\nAccording to the present invention, there is provided a device manufacturing method comprising the steps of:\nexposing a photosensitive agent on a substrate by using an exposure apparatus; and\ndeveloping the photosensitive agent on the substrate,\nthe exposure apparatus having\nan illumination system for illuminating a master,\na master stage apparatus which moves while supporting the master,\na projection optical system for projecting a pattern of the master onto a substrate, and\na substrate stage apparatus which moves while supporting the substrate,\nat least one of the master and substrate stage apparatuses having\na stage,\na member fixed to the stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including the stage and the member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined based on a relative position relationship between the barycentric position of the structure and the at least three spring members so as to set at least one spring constant different from other spring constants.\nAccording to the present invention, there is provided a device manufacturing method comprising the steps of:\nexposing a photosensitive agent on a substrate by using an exposure apparatus; and\ndeveloping the photosensitive agent on the substrate,\nthe exposure apparatus having\nan illumination system for illuminating a master,\na master stage apparatus which moves while supporting the master,\na projection optical system for projecting a pattern of the master onto a substrate, and\na substrate stage apparatus which moves while supporting the substrate,\nat least one of the master and substrate stage apparatuses having\na stage,\na member fixed to the stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including the stage and the member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined to make application points of restoring forces of the at least three spring members coincide with the barycentric position of the structure when the stage is statically balanced.\nAccording to the present invention, there is provided a device manufacturing method comprising the steps of:\nexposing a photosensitive agent on a substrate by using an exposure apparatus; and\ndeveloping the photosensitive agent on the substrate,\nthe exposure apparatus having\na projection optical system for projecting an image onto a substrate, and\na substrate stage apparatus which moves while supporting the substrate,\nthe substrate stage apparatus having\na stage,\na member fixed to the stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including the stage and the member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined based on a relative position relationship between the barycentric position of the structure and the at least three spring members so as to set at least one spring constant different from other spring constants.\nAccording to the present invention, there is provided a device manufacturing method comprising the steps of:\nexposing a photosensitive agent on a substrate by using an exposure apparatus; and\ndeveloping the photosensitive agent on the substrate,\nthe exposure apparatus having\na projection optical system for projecting an image onto a substrate, and\na substrate stage apparatus which moves while supporting the substrate,\nthe substrate stage apparatus having\na stage,\na member fixed to the stage, and\na weight compensation mechanism having at least three spring members for compensating for a weight of a structure including the stage and the member, a barycentric position of the structure deviating from a centroid of a polygon formed by support points of the at least three spring members,\nwherein spring constants of the at least three spring members are determined to make application points of restoring forces of the at least three spring members coincide with the barycentric position of the structure when the stage is statically balanced.\nOther features and advantages of the present invention will be apparent from the following description taken in conjunction with the accompanying drawings, in which like reference characters designate the same or similar parts throughout the figures thereof."}
{"text": "Zeolites have been used extensively to catalyze a number of chemical reactions in refinery and petrochemical reactions, and catalysis, adsorption, separation, and chromatography. For example, both synthetic and natural zeolites and their use in promoting certain reactions, including conversion of methanol to olefins (MTO reactions) and the selective catalytic reduction (SCR) of nitrogen oxides with a reductant such as ammonia, urea or a hydrocarbon in the presence of oxygen, are well known in the art. Zeolites are crystalline materials having rather uniform pore sizes which, depending upon the type of zeolite and the type and amount of cations included in the zeolite lattice, range from about 3 to 10 Angstroms in diameter. Zeolites having 8-ring pore openings and double-six ring secondary building units, particularly those having cage-like structures, have recently been noted to be of interest for use as SCR catalysts. A specific type of zeolite having these properties is chabazite (CHA), which is a small pore zeolite with 8 member-ring pore openings (˜3.8 Angstroms) accessible through its 3-dimensional porosity. A cage like structure results from the connection of double six-ring building units by 4 rings.\nCatalysts employed in the SCR process ideally should be able to retain good catalytic activity over the wide range of temperature conditions of use, for example, 200° C. to 600° C. or higher, under hydrothermal conditions. Hydrothermal conditions are often encountered in practice, such as during the regeneration of a soot filter, a component of the exhaust gas treatment system used for the removal of particles.\nMetal-promoted zeolite catalysts including, among others, iron-promoted and copper-promoted zeolite catalysts, for the selective catalytic reduction of nitrogen oxides with ammonia are known. Iron-promoted zeolite beta (U.S. Pat. No. 4,961,917) has been an effective commercial catalyst for the selective reduction of nitrogen oxides with ammonia. Unfortunately, it has been found that under harsh hydrothermal conditions, for example, as exhibited during the regeneration of a soot filter with temperatures locally exceeding 700° C., the activity of many metal-promoted zeolites begins to decline. This decline is often attributed to dealumination of the zeolite and the consequent loss of metal-containing active centers within the zeolite.\nThe synthesis of a zeolite varies according to structure type of the zeolite, but usually, zeolites are synthesized using a structure directing agent (sometimes referred to as a template, e.g., an organic template) together with sources of silica and alumina. The structure directing agent can be in the form of an organic, e.g., tetraethylammonium hydroxide (TEAOH), or inorganic cation, e.g., Na+ or K+. During crystallization, the tetrahedral silica-alumina units organize around the structure directing agent to form the desired framework, and the structure directing agent is often embedded within the pore structure of the zeolite crystals.\nMetal-promoted, particularly copper-promoted aluminosilicate zeolites having the CHA structure type and a silica to alumina molar ratio greater than 2, particularly those having a silica to alumina ratio greater than or equal to 5, 10, or 15 and less than about 1000, 500, 250, 100 and 50 have recently solicited a high degree of interest as catalysts for the selective catalytic reduction of oxides of nitrogen in lean burning engines using nitrogenous reductants. This is because of the wide temperature window coupled with the excellent hydrothermal durability of these materials, as described in U.S. Pat. No. 7,601,662. Prior to the discovery of metal-promoted zeolites, as described in U.S. Pat. No. 7,601,662, while the literature had indicated that a large number of metal-promoted zeolites had been proposed in the patent and scientific literature for use as SCR catalysts, each of the proposed materials suffered from one or both of the following defects: (1) poor conversion of oxides of nitrogen at low temperatures, for example 350° C. and lower; and (2) poor hydrothermal stability marked by a significant decline in catalytic activity in the conversion of oxides of nitrogen by SCR. Thus, the invention described in U.S. Pat. No. 7,601,662 addressed a compelling, unsolved need to provide a material that would provide conversion of oxides of nitrogen at low temperatures and retention of SCR catalytic activity after hydrothermal aging at temperatures in excess of 650° C.\nIn view of the considerable interest in molecular sieves having 8-ring pore openings and double-six ring secondary building units, particularly those having the CHA structure type, there is an on-going desire to improve the process of preparing such zeolites."}
{"text": "Discovery protocols are known as a technique for automatically acquiring service information of a network device connected on the same LAN. Discovery protocols are found in various forms. Examples of discovery protocols include UPnP (Universal Plug and Play), proposed by Microsoft, and Bonjour, proposed by Apple. There are various other ones, such as WS-Discovery and SLP (Service Location Protocol).\n“Service information” refers to, for example, information (e.g., a service name) for identifying a service such as printing or file storing provided by a network device, and information required for connection, such as the network device name. A network device can easily implement connection with another network device by acquiring the service information of the other device.\nIn a discovery protocol, each device periodically broadcasts its service information over a network or inquires for the service information of a device at the other end, thereby managing a list of services on a LAN. In the case where an inquiry signal is sent to a device with service(s) that has once been found through a search and no response to the inquiry signal is detected, at the time of a service inquiry (search), it is determined that the device has left (disappeared from) the network and that provision of the service(s) in question by that device has stopped, and the connection with the service-providing-device at the other end is then disconnected.\nDisconnecting a connection with a service-providing device is performed in a recognizable manner when the service-providing device has disappeared. The above-described service discovery is intended for a wired network. Accordingly, if there is no response to a service inquiry, it can be determined that the device at the other end has disappeared.\nThis causes a problem specific to wireless communication environments. Examples of such a problem include a temporary deterioration in radio wave conditions and various problems associated with movement of a device. If this kind of problem occurs, throughput decreases due to a temporary deterioration in radio wave condition at a wireless layer. However, connection with a device at the other end is still maintained, and a service-providing device is correctly performing processing.\nIn spite of this, if a service inquiry is sent, no response is obtained, at least temporarily, from the service-providing device. For this reason, a higher-level application which performs service discovery mistakenly determines that provision of a service has stopped and disconnects the lower-level (wireless layer) communication."}
{"text": "The market for portable devices, for example, mobile phones, smart watches, tablets, etc., is expanding with many more features and applications. As the number of applications on these devices increases, there also is an increasing demand to run multiple applications concurrently. More features and applications call for microprocessors to have high performance, but with low power consumption. Multithreading can contribute to high performance in this new realm of application. Keeping the power consumption for the microprocessor and related cores and integrated circuit chips near a minimum, given a set of performance requirements, is desirable, especially in portable device products.\nMultithreading is the ability to pursue two or more threads of control in parallel within a microprocessor pipeline. Multithreading is motivated by low utilization of the hardware resource in a microprocessor. In comparison, multi-core is fairly wasteful of the hardware resource. Multithreading can, in general, provide the same performance as multicore without duplicating of resources.\nMultithreading can be used in an effort to increase the utilization of microprocessor hardware and improve system performance. Multithreading is a process by which two or more independent programs, each called a “thread,” interleave execution in the same processor, which is not a simple problem. Each program or thread has its own register file, and context switching to another program or thread requires saving and restoring of data from a register file to a memory. This process can consume much time and power. These and other problems confront attempts in the art to provide efficient multithreading processors and methods."}
{"text": "1. Field of the Invention\nThe invention relates to an apparatus for extending the lip plate of a dock leveler for movement of the lip plate onto the bed of a truck or other vehicle located next to the dock.\n2. Description of the Prior Art\nDock levelers generally employ a base member attached to the side of a dock, a bridge plate which is pivotally attached to the edge of the base member and a lip plate which is pivotally attached to the other edge of the bridge plate. The known manually operated dock levelers in use today place a large amount of strain on the operator when lifting the plates for proper positioning. Spring mechanisms are employeed on dock levelers to aid in lifting the plates; however, they do not provide,, the assistance needed.\nU.S. Pat. Nos. 3,974,537 and 4,065,824 are directed to different types of mechanisms for assisting the operator in lifting the plates of a dock leveler for proper positioning."}
{"text": "1. Field of the Invention\nThe invention relates to chemical vapor deposition (CVD) systems and, more particularly, to temperature control apparatus, and a concomitant method, for controlling the temperature of the walls of a CVD reaction chamber.\n2. Description of the Prior Art\nIn a conventional chemical vapor deposition (CVD) system, a reaction chamber envelopes a volume wherein the deposition process occurs. The chamber typically is fabricated of transparent quartz. As such, the reaction chamber is said to have quartz windows (accessible surfaces of the chamber) through which the temperature of the apparatus within the chamber can be optically monitored using optical pyrometers. In general, since the windows are integral with the walls of the chamber, the temperature of these windows is indicative of the chamber wall temperature. The temperature of the windows and the walls of the chamber vary depending upon the temperature of the workpiece within the chamber, the process presently being accomplished, the gases involved in the current process, and the run time of the process.\nMore specifically, it is desirable that during a deposition process, the quartz wall temperature be controlled within a narrow temperature range to minimize deposits on the chamber wall. Furthermore, it is desirable that, during cleaning steps, the quartz wall temperature be maintained at a different temperature to maximize the etch rate of any film that may have been deposited on the chamber walls. Such temperature maintenance minimizes particulate contamination problems.\nTypically, the temperature of the walls of the chamber is maintained at a nominal value by a flow of air proximate the outer surface of the chamber. This air flow is fixed for each process performed within the chamber. Thus, the temperature of the chamber walls varies with conditions within the chamber, e.g., the wall temperature differs depending on whether the process in the chamber is depositing material or etching material from a workpiece.\nThrough empirical study, it has been shown that the processes performed within the chamber can be optimized when the chamber wall temperature is maintained at an optimum temperature for a particular process. However, prior art CVD systems contemplate maintaining the air flow past the chamber at a fixed flow rate. As such, these systems exhibit significant chamber wall temperature variations which preclude a high degree of temperature optimization.\nTherefore, a need exists in the art for closed loop apparatus, and a concomitant method, for controlling the air flow proximate a reaction chamber such that the temperature of the chamber walls can be readily optimized for each of the processes accomplished within the chamber."}
{"text": "The scoring of test answer sheets involves complex problems. These test answer sheets typically include a series of response positions such as, for example, \"bubbles,\" ovals, or rectangles. A person taking a test would, for example, darken in an appropriate oval with a pencil to answer a multiple choice question. These test answer sheets may also include handwritten answers, such as essay or short answer questions. Systems for scanning and scoring the bubbles on such answer sheets are known in the art. Increased difficulties are encountered, however, when such answer sheets either include other types of answers, such as handwritten answers, or cannot be machine graded. For example, if the student has failed to include his or her name on the test answer sheet, the system may be unable to machine score the test answer.\nThe goals in scoring test answers that cannot be machine scored include efficiency and consistency. These test answer sheets are typically scored by test resolvers either by manually scoring the physical test answer sheet or scoring an electronic representation of the test answer sheet on a computer. Ideally, the scores provided by the various test resolvers for a particular test question should be consistent, since the scores are used in comparing performance of the students against one another. In addition, a test resolver should ideally work efficiently so as to maintain consistently high scoring rates. The test resolver should not have such a high scoring rate that the consistency or quality of scoring significantly declines; likewise, the test resolver should not have such a low scoring rate that the too few answer sheets are being scored. This manual scoring of test answer sheets, however, makes it difficult to monitor the consistency of scoring among the various test resolvers.\nIn many situations, test resolvers actually travel to a particular location so that all test resolvers may simultaneously score test answer sheets. Requiring the test resolvers to travel to a given location is inconvenient for the resolvers and expensive for those who administer the tests. Furthermore, tracking the performance of test resolvers against both their own performance and the performance of other resolvers can be very difficult with a manual scoring environment.\nThe process of resolving test questions is currently done manually, and this presents problems. A resolver is manually presented with the actual test answer sheets for scoring. This process is relatively inefficient, since the resolvers must score the answer sheets one at a time and in the order in which they are presented. Also, manual scoring systems do not have the capability to efficiently gather and categorize the test answers for subsequent analysis. Therefore, with a manual system it is very difficult to determine how teaching methods should be changed to decrease, for example, the number of incorrect answers.\nA need thus exists for a system that promotes and achieves consistency and efficiency in scoring or resolving of tests."}
{"text": "1. Technical Field\nThe present invention relates to wafer-to-wafer alignments, and more specifically, to wafer-to-wafer alignments using capacitive coupling structures.\n2. Related Art\nIn the prior art, two wafers containing devices can be aligned and then bonded together so as to, among other purposes, double the device density. As a result, there is always a need for a structure (and methods for operating the same) that helps achieve good wafer-to-wafer alignments for subsequent wafer bonding process."}
{"text": "Electromagnetic sources used in geophysical surveying generally produce unfocused regions of subsurface excitation. Unfocused excitations result both from the diffusive nature of low frequency electromagnetic wave propagation in the earth and from the use of relatively simple source configurations in the surveys. Unfocused excitations are a severe disadvantage when a target-oriented survey is needed, such as when it is desirable to probe the electrical properties of a prospective formation in a location where a complex geologic background may be present.\nTo eliminate the disadvantages of unfocused surveys, a method of source focusing is disclosed in U.S. patent application Ser. No. 09/656,191, filed Sep. 6, 2000 (hereinafter referred to as “Srnka”). The method, referred to as the remote reservoir resistivity mapping method (“R3M”) uses a spatially extended grounded electromagnetic source array that is positioned on the earth's surface, or near the seafloor, to focus electromagnetic energy on the subsurface target. The source array consists of two concentric electrode rings whose radii are optimally determined from theoretical modeling to focus the subsurface excitation at the approximate depth of the subsurface target. A limitation of the Srnka method, however, is that for typical reservoir depths, which range from 500 meters to greater than 5000 meters, the total length of electrode wire that is required in the concentric electrode rings ranges from about 23.5 km to greater than 235 km.\nThese extreme lengths of electrode wire raise numerous operational challenges. For example, difficulties in handling and deploying the electrode wires, and problems in ensuring accurate positioning of the wires in the desired electrode ring locations may be involved, and there is the need to avoid obstacles in the area of the desired ring locations, and the requirement that the electrodes be uniformly grounded along their lengths. Other problems influence the economics of the survey, including the large cost of the electrode wire, some of which is likely to be damaged or lost in each survey, and the large labor and logistics costs associated with mobilization, demobilization, deployment, and recovery of the wire.\nFor these reasons, Srnka also disclosed an alternative method that reduces the required length of electrode wire by substituting a set of equally spaced radial electric bipoles, each having a length equal to the difference in the lengths of the radii of the two concentric rings. FIG. 1 depicts both the R3M concentric ring method and this radial bipole alternative. Concentric electrode rings 4 and 5 lie on surface 1 of earth 2 and have radii a and b, respectively, and a center 7 generally above reservoir 3. In the radial bipole alternative method these concentric rings are replaced by radial electrode bipole sources 6. Each bipole source 6 has a length L=b−a. Srnka discloses that a minimum of six radial bipole sources 6 is preferable. A signal equivalent to the signal that a receiver array would receive from the concentric ring source array is produced mathematically by summing the signals received by the receiver array from each of the bipole sources 6 (For convenience, FIG. 1 does not depict the location of the receiver array). This alternative method, which allows the simulation of a concentric ring source array, is referred to as the virtual concentric ring source method, or as the radial bipole method, and substantially reduces the required length of electrode wire as compared to the R3M physical concentric ring source method.\nThe virtual concentric ring source method nevertheless uses relatively long lengths of electrode wire and inherently retains many of the operational challenges of the R3M physical concentric ring source method. In addition, both methods have other inherent disadvantages, in particular for certain aspects of the subsurface resistivity-imaging problem. Specifically, both methods produce maximum responses near a pre-specified subsurface excitation focus depth (in other words, at the approximate depth of the target), but provide little differential sensitivity to features within an appreciable vertical range of that depth. More specifically, the ability to image separately features above or below the target a distance of less than five percent of the target's depth below the surface is limited. This limitation greatly hinders the imaging of stacked reservoirs. In addition, amplifier saturation effects can influence data quality and hinder analysis of such data. A third issue is that both methods produce primarily vertical subsurface currents in the region of the target. These currents are not adequate for imaging reservoir electrical macro-anisotropy (vertical electrical isotropy), such as exists in many reservoirs due to the presence of shale interbeds. Imaging this anisotropy is essential for estimating reservoir net and gross volumes, and for determining hydrocarbon pore volumes.\nAccordingly, a method is desired which retains the advantageous electromagnetic responses of the R3M physical concentric ring source and radial bipole methods, but which reduces or eliminates the disadvantages of those methods. The present invention satisfies that desire."}
{"text": "Several types of modern equipment require an electrical power supply that is capable of producing high energy pulses at a very high pulse repetition rate. Examples of these include, but are not necessarily limited to, certain laser systems, and more modernly, extreme ultraviolet (EAV) sources for advanced photolithography techniques.\nHeretofore, resonant circuits that are configured to charge a capacitor have been employed when power at a high pulse repetition rate has been required. These resonant circuits typically include an inductor in addition to the capacitor. When a power source (e.g. battery) is connected to the resonant circuit and activated, a voltage appears across the capacitor that increases from zero to a maximum value during a time period, t, defined by the LC circuit. After the voltage across the capacitor reaches the maximum value, a load (e.g. UAV laser) can be connected to the resonant circuit to discharge the capacitor and energize the load. This sequence of charging and discharging the capacitor can be repeated, as desired, to drive the load with a train of substantially constant energy pulses.\nOne problem associated with the above-described resonant charge circuit is the regulation of the maximum voltage across the capacitor. Specifically, any residual voltage on the capacitor (i.e. voltage remaining between pulses) will affect the maximum voltage on the capacitor at the end of a charge transfer. In attempts to overcome this difficulty, dequeing techniques have been developed and used. In general, these dequeing techniques have employed a voltage probe to monitor the voltage across the capacitor. When the desired voltage is reached, a switch is used to divert current from the resonant circuit, and as a result, stop all charge transfer to the capacitor.\nBy way of example, FIG. 1 shows a typical, prior art resonant circuit (generally designated 10) that uses a dequeing technique to charge a capacitor 12 at a high pulse repetition rate. As shown, the resonant circuit 10 includes an inductor 14, a power source 16, diodes 18, 19, 20 and two transistor switches 22, 24 (note: the circuit 10 also requires a voltage probe and control circuit that are not shown in FIG. 1). Operation of the resonant circuit 10 begins by closing switch 22 at time t=0. With switch 22 closed, the capacitor 12 is resonantly charged with current that passes through the diode 20 and inductor 14. Once the voltage probe indicates that a desired voltage across capacitor 12 has been reached, the control circuit quickly closes switch 24. With switch 24 closed, all remaining current in the circuit 10 is routed through the circuit branch 21 having switch 22, switch 24 and the inductor 14. With the circuit 10 in this configuration, all current flow to the capacitor 12 is stopped. Switch 22 is then opened, diverting current from branch 21 through a circuit branch 25 having the power source 16, diode 18, inductor 14 and switch 24. This allows the energy in the inductor 14 to be returned to the source 16 and recovered. The charge across the capacitor 12 can then be maintained until required by a load (not shown). Once the capacitor 12 has been discharged, switch 24 is then opened, configuring the circuit 10 to generate the next pulse.\nIn a typical setup of the prior art resonant circuit 10 shown in FIG. 1, the desired voltage across the capacitor 12 is selected to be less than the peak voltage generated by the resonant circuit 10, which in turn, is typically about twice the voltage of the source 16. For the circuit 10, the switches 22, 24 are preferably constructed of either MOSFET's or IGBT's, which unfortunately, have limited voltage ratings. Specifically, operational charging voltages for the circuit 10, as shown, have been generally limited to a maximum voltage that is below about 2 kV.\nIn light of the above, it is an object of the present invention to provide a power supply that is capable of producing pulses at a very high pulse repetition rate and that is operable at relatively high voltages (i.e. greater than about 2 kV). It is yet another object of the present invention to provide a system for charging a capacitor at a high voltage and high pulse repetition rate while accurately regulating the maximum voltage across the capacitor. Yet another object of the present invention is to provide a high pulse rate, pulsed power system which is easy to use, relatively simple to implement, and comparatively cost effective."}
{"text": "The present invention generally relates to tape cassettes, and more particularly to a tape cassette having a construction such that one reel of a supply side reel and a take-up side reel, for example, the supply side reel, is fitted over a reel driving shaft of a recording and/or reproducing apparatus and rotated at a center part thereof, and the other reel, for example, the take-up side reel, is transmitted of rotation at a part other than a center part thereof and rotated.\nPresently, video signal recording and/or reproducing apparatuses using tape cassettes are reduced to practical use on the world-wide basis. As types of these recording and/or reproducing apparatuses using tape cassettes, two or three standardized types of apparatuses presently exist world-widely. There is no interchangeability between apparatuses of different standards, however, interchangeability exists between apparatuses adopting the same standard. Accordingly, a tape cassette recorded by one recording and/or reproducing apparatus can be reproduced by another recording and/or reproducing apparatus of the same standard. Hence, in order to ensure interchangeability between different recording and/or reproducing apparatuses manufactured by different manufacturers so as to perform recording and/or reproduction, a standardization is set with respect to formats of the tape cassette and the recording and/or reproducing apparatus. That is, the above formats are standardized for each standard, and a standard type tape cassette and a standard type recording and/or reproducing apparatus are respectively provided and employed for each standard.\nRecently, a problem of much importance involves the realization in reducing the size of the recording and/or reproducing apparatus main body. This is to realize a more compact portable type recording and/or reproducing apparatus, and, for example, to realize such an equipment that a recording apparatus is unitarily built into a television camera.\nIn this regard, some attempts have been made to realize a compact type recording and/or reproducing apparatus using a compact or miniature type tape cassette, which uses a tape pattern and format completely different from those of the standard type recording and/or reproducing apparatus using the standard type tape cassette which are already and widely accepted in the market. However, the above compact type recording and/or reproducing apparatus attempted for realization uses a format completely different from that of the standard type recording and/or reproducing apparatus, and interchangeability does not exist therebetween. Therefore, in this compact type recording and/or reproducing apparatus, there is a disadvantage in that a tape cassette recorded by the compact type recording and/or reproducing apparatus cannot be reproduced by the standard type recording and/or reproducing apparatus. This is a great inconvenience and disadvantage to the owners of the compact type recording and/or reproducing apparatus.\nAnother attempt has also been made to realize a compact type portable recording and/or reproducing apparatus which performs recording and/or reproduction with the same tape pattern and format as the standard type recording and/or reproducing apparatus, by using a tape cassette whose size is slightly reduced by reducing the tape quantity and the diameter of the reels from those of the standard type tape cassette. In this system, a tape cassette recorded by the portable recording and/or reproducing apparatus can be reproduced as it is by the standard type recording and/or reproducing apparatus.\nHowever, in the tape cassette used in the above portable recording and/or reproducing apparatus, the distance between a supply side reel and a take-up side reel is set equal to the distance between the supply side and take-up side reels of the standard type tape cassette, so that the tape cassette used for the portable recording and/or reproducing apparatus can be loaded into the standard recording and/or reproducing apparatus. Hence, even when the tape quantity is reduced in order to reduce the diameter of the reels, there is a limit in reducing the diameter of these reels. Therefore, in this system, there was a disadvantage in that the size of the tape cassette as a whole could not be reduced significantly, and the same is true to the recording and/or reproducing apparatus.\nFurthermore, another system can be considered in which the tape quantity is reduced to reduce the diameter of the reels, and further, the miniature type tape cassette is constructed by reducing the distance between the supply side and take-up side reels. In this system, as a modification of the standard type recording and/or reproducing apparatus, the recording and/or reproducing apparatus can be constructed so that a reel disc for driving the take-up side reel is movable. In this modification of the standard type recording and/or reproducing apparatus, the take-up side reel disc is at a normal position when loaded with the standard type tape cassette, and the take-up side reel disc is moved to a position closer to the supply side reel disc when loaded with the miniature type tape cassette.\nHowever, even in this system, for example, there is a disadvantage in that the miniature type tape cassette recorded by the compact type recording and/or reproducing apparatus cannot be reproduced by the standard type recording and/or reproducing apparatus which is presently in wide use. Moreover, there is a disadvantage in that it is extremely difficult to realize a mechanism for moving the reel disc in the manner described above. Furthermore, the recording and/or reproducing apparatus using the tape cassette generally has a mechanism for drawing out the tape from within the tape cassette to load the tape onto a predetermined tape path within the recording and/or reproducing apparatus. Therefore, it is also exceedingly difficult to construct the above mechanism for pulling out the tape, so that interchangeability exists with respect to the above compact or miniature type tape cassette and the standard type tape cassette, and the realization of such a mechanism is virtually impossible."}
{"text": "Payment devices, such as credit/debit cards or contactless devices (e.g., “PAY WAVE”™) are used in various financial transactions including credit purchases, debit purchases, cash advances, and payroll transactions. During the financial transaction, authorization messages that include the financial details about the transaction are transmitted from a point-of-sale (POS) terminal located at a merchant to an acquirer (e.g., a financial institution). The acquirer obtains from the financial institution that issued the payment device (“the issuer”) information regarding the credit worthiness of the cardholder, and whether the account has sufficient funds or credit to perform the transaction.\nThe acquirer then forwards the financial transaction information to a payment processor, such as Visa, MasterCard, or American Express, which then performs the remainder of the financial transaction along with the issuer (e.g., commercial bank through which the payment device was issued). The only data that are transmitted in these conventional payment device financial transactions identify the amount of the transaction, the account (e.g., card number) performing the transaction and possibly the location where the transaction is occurring. Aside from these details, little is known about the financial transaction.\nAn improved system and method of performing a financial transaction is desirable."}
{"text": "Video codec consists of an encoder that transforms the images of the input video into a compressed representation suited for storage/transmission and a decoder that can uncompress the compressed video representation back into a viewable form. Image sequences can be compressed either as sequences of still pictures coded with spatial prediction means or inter pictures coded with spatial and temporal prediction means. Image sequences with random access and support for editing individual pictures have been traditionally enabled by representing the sequence as a series of independently coded intra pictures. Such formats include, for example, Motion JPEG, animated GIF and the Intra profiles of H.264.\nIf a sequence of images is represented as a series of still pictures, the coding efficiency is typically poor and the file size requirement for a high resolution sequence can become massive. In order to access a predicted picture, the decoder needs to conclude which reference pictures are required to be decoded first. The delay may be significantly reduced, if only those pictures that are used as reference are decoded rather than decoding all pictures starting from the previous IDR picture or alike. However, in the known methods, random access to any sample may require decoding all reference pictures starting from the previous IDR picture or alike, as concluding which subset of reference pictures are required to be decoded may not be easily possible.\nTherefore there is a need for an improved method for identifying pictures that are required to be decoded in order to decode a desired random-accessed picture correctly."}
{"text": "This is a continuation of application Ser. No. 10/758,654 filed Jan. 15, 2004 now U.S. Pat. No. 6,960,140.\nThis invention relates generally to golf equipment and, in particular, to a golf putter head with a visual alignment aid and an increased moment of inertia.\nRecent developments in golf equipment have resulted in golf putter heads with high moments of inertia. For example, U.S. Pat. No. 5,482,281 to D. W. Anderson discloses a putter head sold under the name DANSER. The Anderson putter head has heel and toe weights mounted on a lower plate-like member. The heel and toe weights and the lower plate-like member are preferably made of heavyweight material such as bronze or steel. An upper shell-like member, preferably made of lightweight material such as plastic or aluminum, is secured to the lower plate-like material to enclose the heel and toe weights. U.S. Pat. No. 5,842,935 to M. J. Nelson discloses a putter head sold under the name NELLI. The Nelson putter head has a horseshoe shaped body formed of high density material such as steel with thickened heel and toe portions. The horseshoe shaped body includes a cavity which receives an insert formed of low density material such as polyurethane. The insert preferably constitutes about 15% of the total weight of the putter head while constituting more than 50% of the total volume of the putter head."}
{"text": "The analysis of cell samples is a valuable tool in the treatment of disease. If a patient complains of an issue with their body, or portion of tissue appears unusual in appearance, it is common practice to obtain a sample of the cells at that site, and to analyze the sample in a lab. If the tissue is hard, such as that found with a hard tumor, a biopsy coring needle or a rasp or a scraper can be used. If the surgical site is softer tissue such as the mucosal or submucosal tissue found within the gastrointestinal tract, a scraper, rasp or cytology brush can be used.\nCytology brushes are well known in the art, and have been available to the medical community for a very long time. Flexible shafts were added to the cytology brushes to increase access to hard to reach portions of the patient. With the advent of endoscopic surgery, cytology brushes were adapted for endoscopic use. As endoscopes advanced, more and more hard-to-reach areas became available for cytological brush sampling. Today, flexible shaft cytology brushes can be used during endoscopic surgery to access difficult-to-reach portions of the patient's anatomy. Tissue samples can now be taken with a cytology brush from within the difficult-to-access biliary tract, which can be of great help in the early diagnosis of diseases of the gall bladder and pancreas.\nConsequently, a significant need exists for an improved cytology brush apparatus that can access difficult-to-reach anatomy, and provide improved cell collection at the surgical site."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of tropospheric scatter radio links and more particularly to a radiating system with angular diversity comprising a main reflector, a subreflector, a transmitting horn and at least two receiving horns.\n2. Description of the Prior Art\nIt is known that to establish microwave radio links beyond the horizon it is possible to use radiating systems which utilize the scattering of electromagnetic waves by the troposphere.\nIt is also known that the troposphere displays irregularities generally considered as bubbles or layers which vary continuously in number, form and position with resulting variation of the refraction index and diffusion angle. When such irregularities are illuminated by a beam of electromagnetic waves from a transmitting antenna they scatter the electromagnetic energy in all directions but predominantly within a cone having as its axis the direction of transmission.\nIt is clear that with such links path attenuation is much higher than that found in links with antennas which remain in a field of mutual visibility since the propagation mechanism is different. In addition, in troposcatter radio links there are met sudden deep fadings of the intensity of the signal received due mainly to random movements of the irregularities of the troposphere.\nDiversity techniques are known which are used to avoid the aformentioned problems with tropospheric propagation, i.e. spatial, frequency and angular diversity. Diversity can also be simple or multiple. In case of multiple diversity suitable combinations of the different diversity techniques have been achieved.\nSpatial diversity consists of transmitting the same signal with two antennas appropriately spaced and directed and in using two other antennas similarly arranged for reception. The basic assumption on which this technique is based is that fadings of signal intensity which appear on the two beams are poorly correlated.\nFrequency diversity differs from spatial diversity in that the signal is radiated on a single beam but with two carriers appropriately spaced as to frequency so as to make intensity fadings of the two signals received uncorrelated.\nAngular diversity consists of radiating electromagnetic power in a single beam and in equipping the receiving antenna with two receiving horns appropriately spaced from each other in such a manner that the single transmitted beam is received in two different directions forming a certain angle called diversity angle and giving rise to two signals as independent as possible from the point of view of tropospheric propagation. It is thus possible to effect in reception a combination of the two signals received such that the combination signal intensity or the signal-to-noise ratio of the combination is always kept sufficiently high.\nIt is also known that with angular diversity systems there is the problem of optimizing the diversity angle which, as mentioned above, depends on the distance between the receiving horns. As the diversity angle increases so does the statistical independence between the intensity fadings which appear on the two received signals, with a resulting system improvement. But antenna gain is simultaneously reduced because of defocusing.\nIt is also known that radiating systems in general and those with angular diversity in particular accomplish the transmitting part and the receiving part on the same antenna and bring about decoupling of the transmitting signals from the receiving signals by using different frequencies or by means of polarizations on the orthogonal planes or with a combination of these decoupling criteria. As concerns polarization, there are radiating systems with single polarization and radiating systems with double polarization.\nRadiating systems with double-polarization angular diversity possess a first horn generally placed in the focus of the antenna parabola used for both transmitting and receiving and a second horn arranged parallel to the first used only for receiving.\nThe drawbacks of systems of this type are due mainly to the complexity of antenna horns. In general they include for effecting decoupling or discrimination between the two orthogonal polarizations many elements which lead to considerable occupied space with the resulting reduction of efficiency of the antenna compared with theoretical efficiency.\nAn example of a known tropospheric radiating systems with single-polarization angular diversity is British Pat. No. 1,178,782 granted Jan. 21, 1970 to the Marconi Company Limited which utilizes a parallel-conductor screen to separate the reception polarization from the transmission polarization, which are orthogonal to each other.\nThe above system described in the aforementioned British patent makes use of an offset paraboloid to permit the beam leaving the transmitting horn placed in the focus of said paraboloid to reach the surface of the antenna, avoiding blocking effects by the receiving horns which are outside the field of illumination.\nThe drawbacks of the angular-diversity radiating system described are due mainly to the fact that in such a system the primary illumination axis forms an offset angle with the orthogonal optical axis at the antenna aperture plane. As is known, offset systems provide performance generally poorer than symmetrical systems and in particular have less efficiency in crossed polarization because as is known for efficiency diminishes as antenna curvature increases, i.e. for smaller focus-to-diameter ratios and especially for geometrical dissymetries of the optical system.\nThe drawbacks mentioned hereinbefore are all the more serious in prior art systems which, as described hereinbefore, use a parallel-conductor screen to separate the two linear polarizations which are othogonal to each other. As a result of less efficiency under crossed polarization, a part of the electromagnetic power of the transmitted beam leaves the antenna with a polarization orthogonal to that which it should have. This part of the power, after reaching the receiving antenna, passes through the parallel-conductor screen and reaches the transmitting horn while it should be reflected from the screen toward the receiving horns.\nAccordingly the primary object of the present invention is to overcome the aforementioned drawbacks of the prior art and provide an angular-diversity radiating system which is symmetrical, permits the use of antenna horns which are easy to fabricate, and has good efficiency under crossed polarization, and permits adjustment of the distance between the receiving horns to optimize the diversity angle."}
{"text": "These teachings relate to optical designs, which provide reduced unwanted background radiation, increased signal throughput, or a combination of both.\nAs detecting elements become smaller, the size of the associated slit elements in imaging spectrometer designs typically become smaller. With this decreased size, the effects of the diffraction of light become more significant, leading to the leakage of unwanted background radiation from outside the geometric rays to pass through to the detecting element, as well as leading to the loss of signal radiation that is diffracted outside of the geometric ray bundle and vignette by the optical elements or stop of the system.\nThis effect can be quite significant, particularly in infrared systems where unwanted background radiation that is not shielded from the detecting element can prematurely fill the electron wells of photonic detecting elements, or mask the true temperature of objects in the scene in the case of pyroelectric detecting elements. In systems where radiometric sensitivity or accuracy is concerned, the impact of these diffraction effects can be intolerable.\nFor example, consider some applications of hyperspectral imaging in which it is desirable to have a spectrometer that is intended to identify targets by their spectral signatures in the infrared portion of the electromagnetic spectrum, but must also be very small size and mass so that the system can be transported in an unmanned aerial vehicles (UAV) or be man-portable. These sensors may typically have detecting elements and associated slit apertures whose dimensions are on the same order as the wavelength of the light they are intended to detect, in which case the effects of diffraction are substantial, and can easily reach diffraction angles of 30 degrees or more. If accurate radiometric measurement of the targets is necessary, then this diffraction can mask or hide the target from detection or even introduce false alarms detections. Additionally, the loss of wanted signal from the target due to diffraction can reduce the sensitivity of the sensor, allowing potential targets and threats to go unidentified.\nThere is a need for optical designs that increasing the overall throughput, decrease the overall background radiation, or a combination thereof."}
{"text": "Many modern communication systems use time division multiplexing. Radios in such multiplexed systems switch rapidly between the receive and transmit modes. Such radios may include oscillators that provide reference waveforms for both the receiver and transmitter portions of the radios. A problem may arise in such radios when an oscillator in the radio produces a signal having a frequency equal to the receive frequency, or to an intermediate frequency of the receiver, because in many cases, the oscillator cannot be turned off sufficiently fast to prevent its output from being received by the receiver when the radio goes into its receive mode. These oscillator signals interfere with the performance of the radio because the receiver receives the oscillator signals instead of desired signals. This phenomenon is known as self-quieting, and is obviously undesirable. One possible solution to the self-quieting problem is to shield the receiver portion of the radio so that the oscillator signals are not received during the receive mode. However, shielding of a receiver requires space and may not be as effective as desired. Thus, a need exists for a method and apparatus for reducing self-quieting that avoids the detriments of shielding."}
{"text": "The need to regulate the gas flow between the barrel and operating system of a firearm has been a concern since the introduction of auto-loading firearms. Gas is generated during the combustion of gun powder present in the cartridges used in modern firearms. This gas expands violently to push the bullet out of the firearm's barrel. These expanding gases are utilized as a means to operate the action of the host firearm. In modern firearms the preferred method of facilitating the function of an auto-loading weapon is as follows. A hole is placed through the barrel, generally on the top. Location of this hole or gas port varies between operating systems. Generally a gas port size is chosen to allow a broad range of ammunition to be utilized while guaranteeing the reliable function of the host firearm. Unfortunately due to varying lengths of barrels, ammunition variance, and other factors it is very difficult to choose a gas port size which universally works under all conditions. A popular way of dealing with these problems is to incorporate an adjustable gas block into the operating system.\nAn adjustable gas block allows for the flow of gas between the gas port in the barrel and the operating system of the firearm to be increased or decreased based on mitigating factors present at the time of use. These systems typically work by utilizing an oversized gas port with means to adjust the flow of gas into the operating system and by venting the unneeded gases from the barrel into the atmosphere thus generating flash and sound. Further, adjustment of the gas system typically requires a special tool and offers no way for the user to index the system and make adjustments due to mitigating circumstances quickly. Designs such as these are well known in the prior art and can be found on the Belgium FAL, Soviet SVD and the Yugoslavian M76 rifle.\nRecent firearm designs such as the FN SCAR rifles have incorporated adjustable gas blocks to be used in conjunction with noise suppressors. Noise suppressors provide a means to redirect, cool and slow the expanding gases generated from the discharge of a firearm so that the resulting flash and sound generated by the firearm is minimized or eliminated. As a result, back pressure is generated forcing more gas into the firearm's operating system. This extra gas, or back pressure increases the firing rate of a weapon during its full auto function, fouls the weapon leading to premature malfunction and to a variety of feeding and extraction problems.\nProblems with existing systems may occur due to variations in cartridge, such as different weights, different powder charges, different bullet jackets, etc., friction in the operating system that may change during the life of the rifle, buffer sprint set changes, and suppressors of different back-pressures. Existing systems may not allow for the rifle operator to compensate for these potential problems.\nNeeds exist for improved systems and methods for gas blocks for self-loading firearms to facilitate user adjustment of gas flow from a barrel into an operating system."}
{"text": "This invention relates to printing apparatus, particularly for highspeed rotary multicolour web printing as by the offset process in the production of newspapers and periodicals though the invention may have application to other types of single or multicolour web or sheet fed printing processes and apparatus.\nIt is known to provide printing apparatus having individual process units each driven mechanically independently by its own electric motor instead of coupling the units together by shafts, gearing etc to be driven in synchronism from a common motor. The individual motors are of a type which can be extremely accurately controlled electrically as to running speed and angular position of their drive output to provide synchronisation of the various units and accurate registration of their respective operations on the web or material being processed. This kind of drive arrangement is commonly referred to as xe2x80x9cshaftless drivexe2x80x9d and many examples are well known in the art such as in our GB 2149149A.\nAlthough shaftless drive is successful in providing accurate and consistent results when the printing apparatus is running at its service speed, problems arise which have not hitherto been solved during changes in operating speed of the apparatus as a whole and, consequently, the linear speed of the web, due mainly to stretching or other distortion of the web consequent upon the speed change. This results in the operations carried out by the different units becoming out of register, for example registration of different colours in multicolour printing is lost giving rise to substantial wastage and loss of production time.\nThe apparatus has to be run up to speed on start-up and has to be run down again to close down giving wastage at the beginning and end of every print run.\nThe object of the invention is to provide printing apparatus having shaftless drive with all its consequent advantages and which is able to maintain accurate synchronisation and registration throughout substantial changes in overall running speed.\nAccording to a first aspect of the invention there is provided drive means for printing apparatus of the kind comprising a plurality of process units each for effecting a respective repeat operation on web or sheet material fed along a feed path defined at least in part by said units; said drive means comprising an electrically powered rotary motor individual to each said unit or to one or more rotary elements within a unit, and control means operating to regulate the running speed and angular phasing of each said motor individually in relation to a predetermined master standard programmed in to the control means: characterised in that the control means further includes speed correction means for varying the relative phase angles of one or more of the motor drive outputs by a preselected factor proportionate to the overall running speed of the apparatus to maintain accurate registration relative to the material of the various operations carried out by said units throughout a substantial range of operating speeds.\nThe invention further resides in a method of operating printing apparatus incorporating shaftless drive including the steps of determining the extent of registration error caused by changes in overall operating speed taking place without any adjustment in relative speed of individual drive motors or changes in relative angular phasing of their drive outputs to provide an algorithm representing register error as a function of said operating speed, and applying that algorithm automatically to vary the relative phase angles of one or more of the motor drive outputs by a preselected factor proportionate to the overall running speed of the apparatus whereby registration is correctly maintained throughout a substantial range of operating speeds."}
{"text": "The present invention relates to introduction of naked DNA and RNA sequences into a vertebrate to achieve controlled expression of a polypeptide. It is useful in gene therapy, vaccination, and any therapeutic situation in which a polypeptide should be administered to cells in vivo.\nCurrent research in gene therapy has focused on xe2x80x9cpermanentxe2x80x9d cures, in which DNA is integrated into the genome of the patient. Viral vectors are presently the most frequently used means for transforming the patient\"\"s cells and introducing DNA into the genome. In an indirect method, viral vectors, carrying new genetic information, are used to infect target cells removed from the body, and these cells are then re-implanted. Direct in vivo gene transfer into postnatal animals has been reported for formulations of DNA encapsulated in liposomes and DNA entrapped in proteoliposomes containing viral envelope receptor proteins (Nicolau et al., Proc. Natl. Acad Sci USA 80:1068-1072 (1983); Kaneda et al., Science 243:375-378 (1989); Mannino et al., Biotechniques 6:682-690 (1988). Positive results have also been described with calcium phosphate co-precipitated DNA (Benvenisty and Reshef Proc. Natl. Acad Sci USA 83:9551-9555 (1986)).\nThe clinical application of gene therapy, as well as the utilization of recombinant retrovirus vectors, has been delayed because of safety considerations. Integration of exogenous DNA into the genome of a cell can cause DNA damage and possible genetic changes in the recipient cell that could predispose to malignancy. A method which avoids these potential problems would be of significant benefit in making gene therapy safe and effective.\nVaccination with immunogenic proteins has eliminated or reduced the incidence of many diseases; however there are major difficulties in using proteins associated with other pathogens and disease states as immunogens. Many protein antigens are not intrinsically immunogenic. More often, they are not effective as vaccines because of the manner in which the immune system operates.\nThe immune system of vertebrates consists of several interacting components. The best characterized and most important parts are the humoral and cellular (cytolytic) branches. Humoral immunity involves antibodies, proteins which are secreted into the body fluids and which directly recognize an antigen. The cellular system, in contrast, relies on special cells which recognize and kill other cells which are producing foreign antigens. This basic functional division reflects two different strategies of immune defense. Humoral immunity is mainly directed at antigens which are exogenous to the animal whereas the cellular system responds to antigens which are actively synthesized within the animal.\nAntibody molecules, the effectors of humoral immunity, are secreted by special B lymphoid cells, B cells, in response to antigen. Antibodies can bind to and inactivate antigen directly (neutralizing antibodies) or activate other cells of the immune system to destroy the antigen.\nCellular immune recognition is mediated by a special class of lymphoid cells, the cytotoxic T cells. These cells do not recognize whole antigens but instead they respond to degraded peptide fragments thereof which appear on the surface of the target cell bound to proteins called class I major histocompatibility complex (MHC) molecules. Essentially all nucleated cells have class I molecules. It is believed that proteins produced within the cell are continually degraded to peptides as part of normal cellular metabolism. These fragments are bound to the MHC molecules and are transported to the cell surface. Thus the cellular immune system is constantly monitoring the spectra of proteins produced in all cells in the body and is poised to eliminate any cells producing foreign antigens.\nVaccination is the process of preparing an animal to respond to an antigen. Vaccination is more complex than immune recognition and involves not only B cells and cytotoxic T cells but other types of lymphoid cells as well. During vaccination, cells which recognize the antigen (B cells or cytotoxic T cells) are clonally expanded. In addition, the population of ancillary cells (helper T cells) specific for the antigen also increase. Vaccination also involves specialized antigen presenting cells which can process the antigen and display it in a form which can stimulate one of the two pathways.\nVaccination has changed little since the time of Louis Pasteur. A foreign antigen is introduced into an animal where it activates specific B cells by binding to surface immunoglobulins. It is also taken up by antigen processing cells, wherein it is degraded, and appears in fragments on the surface of these cells bound to Class II MHC molecules. Peptides bound to class II molecules are capable of stimulating the helper class of T cells. Both helper T cells and activated B cells are required to produce active humoral immunization. Cellular immunity is thought to be stimulated by a similar but poorly understood mechanism.\nThus two different and distinct pathways of antigen processing produce exogenous antigens bound to class II MHC molecules where they can stimulate T helper cells, as well as endogenous proteins degraded and bound to class I MHC molecules and recognized by the cytotoxic class of T cells.\nThere is little or no difference in the distribution of MHC molecules. Essentially all nucleated cells express class I molecules whereas class II MHC proteins are restricted to some few types of lymphoid cells.\nNormal vaccination schemes will always produce a humoral immune response. They may also provide cytotoxic immunity. The humoral system protects a vaccinated individual from subsequent challenge from a pathogen and can prevent the spread of an intracellular infection if the pathogen goes through an extracellular phase during its life cycle; however, it can do relatively little to eliminate intracellular pathogens. Cytotoxic immunity complements the humoral system by eliminating the infected cells. Thus effective vaccination should activate both types of immunity.\nA cytotoxic T cell response is necessary to remove intracellular pathogens such as viruses as well as malignant cells. It has proven difficult to present an exogenously administered antigen in adequate concentrations in conjunction with Class I molecules to assure an adequate response. This has severely hindered the development of vaccines against tumor-specific antigens (e.g., on breast or colon cancer cells), and against weakly immunogenic viral proteins (e.g., HIV, Herpes, non-A, non-B hepatitis, CMV and EBV).\nIt would be desirable to provide a cellular immune response alone in immunizing against agents such as viruses for which antibodies have been shown to enhance infectivity. It would also be useful to provide such a response against both chronic and latent viral infections and against malignant cells.\nThe use of synthetic peptide vaccines does not solve these problems because either the peptides do not readily associate with histocompatibility molecules, have a short serum half-life, are rapidly proteolyzed, or do not specifically localize to antigen-presenting monocytes and macrophages. At best, all exogenously administered antigens must compete with the universe of self-proteins for binding to antigen-presenting macrophages.\nMajor efforts have been mounted to elicit immune responses to poorly immunogenic viral proteins from the herpes viruses, non-A, non-B hepatitis, HIV, and the like. These pathogens are difficult and hazardous to propagate in vitro. As mentioned above, synthetic peptide vaccines corresponding to viral-encoded proteins have been made, but have severe pitfalls. Attempts have also been made to use vaccinia virus vectors to express proteins from other viruses. However, the results have been disappointing, since (a) recombinant vaccinia viruses may be rapidly eliminated from the circulation in already immune individuals, and (b) the administration of complex viral antigens, may induce a phenomenon known as xe2x80x9cantigenic competition,xe2x80x9d in which weakly immunogenic portions of the virus fail to elicit an immune response because they are out-competed by other more potent regions of the administered antigen.\nAnother major problem with protein or peptide vaccines is anaphylactic reaction which can occur when injections of antigen are repeated in efforts to produce a potent immune response. In this phenomenon, IgE antibodies formed in response to the antigen cause severe and sometimes fatal allergic reactions.\nAccordingly, there is a need for a method for invoking a safe and effective immune response to this type of protein or polypeptide. Moreover, there is a great need for a method that will associate these antigens with Class I histocompatibility antigens on the cell surface to elicit a cytotoxic T cell response, avoid anaphylaxis and proteolysis of the material in the serum, and facilitate localization of the material to monocytes and macrophages.\nA large number of disease states can benefit from the administration of therapeutic peptides. Such peptides include lymphokines, such as interleukin-2, tumor necrosis factor, and the interferons; growth factors, such as nerve growth factor, epidermal growth factor, and human growth hormone; tissue plasminogen activator; factor VIII:C; granulocyte-macrophage colony-stimulating factor; erythropoietin; insulin; calcitonin; thymidine kinase; and the like. Moreover, selective delivery of toxic peptides (such as ricin, diphtheria toxin, or cobra venom factor) to diseased or neoplastic cells can have major therapeutic benefits. Current peptide delivery systems suffer from significant problems, including the inability to effectively incorporate functional cell surface receptors onto cell membranes, and the necessity of systemically administering large quantities of the peptide (with resultant undesirable systemic side effects) in order to deliver a therapeutic amount of the peptide into or onto the target cell.\nThese above-described problems associated with gene therapy, immunization, and delivery of therapeutic peptides to cells are addressed by the present invention."}
{"text": "For the owners of sporting rifles there exists a recurrent necessity to test-fire their guns, e.g. for adjusting a telescopic sight or a diopter sight or for determining the accuracy after a change of ammunition. Testing of a gun requires a certain number of shots to be fired and cannot, however, be carried out without interruption, since the barrel of the gun under test heats up rather quickly.\nThis requires the sporting gun to be stood aside for cooling or, in order to avoid it being damaged, to be retained in a gun stand, inasmuch as otherwise the gun barrel might be harmfully affected by overheating. Cooling of the gun barrel requires considerable time, however, so that test firing of a sporting gun is a rather time-consuming operation.\nIt is thus an object of the invention to provide means enabling the period of time required for cooling a shot-heated gun barrel to be considerably shortened.\nFor attaining this object, the invention provides that the frame comprises a conduit extending at an inclined position relative to the ground for receiving at least the barrel of the supported sporting gun and being adapted to have a gaseous cooling medium conducted therethrough.\nIn this manner it is possible to considerably shorten the time required for cooling of a shot-heated gun barrel, so that test-firing of a sporting gun requires only short intervals.\nIn greater detail the invention provides that a blower driven by an electric motor is connected to the lower end of the conduit for feeding a gaseous cooling medium thereto.\nUse is made of a high-performance blower. The inclined position of the conduit receiving the heated gun barrel results in the cooling air flow being directed obliquely upwards, so that the possibly warm air stream will not whirl up any loose particles from the ground or molest any persons staying in the shooting stand.\nThe invention further provides that the conduit gradually widens towards its upper end and has its inner side insulated by means of a heat-resistant, resiliently flexible material.\nThis provision ensures that the gun barrel to be cooled may be introduced into the conduit e.g. with a telescopic sight attached thereto, and that the cross-section of the conduit has only to be dimensioned such that the sporting gun, or the barrel, respectively, is enclosed on all sides by an air stream just adequate to carry off the heat. The lining or insulation of the conduit serves to protect the gun barrel against damage in case of it being introduced less carefully. At the same time the turbulence of the air stream along the insulation or lining of the conduit prevents heat transmission thereto.\nIn an advantageous embodiment of the support apparatus the invention provides that the conduit with its blower is supported on a tubular frame member extending at an inclined position relative to the ground and carrying the support members for the respective sporting gun, said frame member comprising a telescopically extensible and lockable tubular profiled rod carrying at its projecting upper end a pair of support members in the form of forks disposed at an angular spacing of 90.degree., one of said forks serving as a support member for a gun stock extending substantially in the direction of the gun barrel, and the other fork serving as a support for a gun stock extending at an angle to the gun barrel.\nThis embodiment enables the support apparatus to be adapted to the length of the sporting gun, since the support member for the gun stock is adjustable with respect to the upper end of the conduit, so that in any case the entire gun barrel extends into the conduit irrespective of the configuration of the remainder of the gun. The forks are disposed at the projecting upper end of the tubular profiled rod in such a manner that the gun barrel is positioned substantially at the longitudinal center of the conduit adjacent the upper end thereof.\nIn order to ensure that the gun barrel to be cooled extends substantially parallel to the longitudinal center axis of the conduit over its entire length the conduit contains a gun barrel centering support comprising a U-shaped yoke supported at the free ends of its legs for pivotal movement about a pivot axis extending transversally of the longitudinal direction of the conduit and extending into the conduit, the center portion of said yoke being bent into a V-shape for centering the gun barrel, with the apex of said V-shape being directed towards said pivot axis. The pivot axis of the centering yoke is disposed at the underside of the tubular frame member in such a position that even the shortest barrel of a sporting gun may be securely centered.\nFor adjusting the centering support, the pivot axis of the centering yoke is provided at least at one end thereof with an adjusting lever adapted for pivotal movement between two limit positions, said pivot axis or said adjusting lever carrying a cam engaging a cam follower of a cam follower limit switch for energizing the blower in the limit position corresponding to the centering position of the yoke. This enables centering of the gun barrel and starting of the blower to be carried out in a single operation.\nFinally the invention provides that the frame is movable along the ground on ground-engaging rollers or, in an open-air embodiment, on rubber-tired cross-country wheels of larger diameter, and is stationarily positionable on rubber pads in the manner of a tripod. The support apparatus may thus be movably supported on the ground-engaging rollers and be brought to rest substantially immovably on the rubber pads."}
{"text": "A syringe of the two-compartment type embodying the structure described above and wherein a bypass is provided through which the liquid can be forced to flow into the compartment containing an agent to be dispersed in that liquid, is described in Europatent publication No. 85 101 508.\nIn this syringe, one plunger part forms a partition between the two compartments while the other plunger part is attached to the stem and, as the stem is advanced, the second plunger part displaces the liquid and the first plunger part past the bypass until the two plunger parts abut one another, whereupon the plunger formed by these parts is displaced by the stem to eject the composition from the needle.\nThe syringe described in the Europatent publication is a so-called disposable syringe which is sterilized and prefilled with the injectable substance in the two-component state.\nSuch two-compartment syringes have been found to be especially effective when the injectable substance is to be in the form of a mixture in a liquid phase or a solution of a component in the liquid which has only a limited effective life in the liquid composition and thus must be fabricated as shortly before use as is possible.\nThe two components are thus provided in separate compartments and remain stable in these compartments for long periods of time.\nSince one compartment is provided axially behind the other and the two substances are only brought into contact with each other immediately prior to injection, it is necessary to provide a means for permitting the liquid from the compartment most distal from the needle to mix with the solid substance, for example, in the compartment proximal to the needle.\nThis is achieved by providing the bypass immediately ahead of the plunger member which delimits the liquid compartment at its side closest to the needle.\nConsequently, as the plunger stem is activated, i.e. displaced into the cylinder or barrel of the syringe, the plunger part most distal from the needle will drive the body of liquid ahead of it and the other plunger part toward the needle until communication is established between the compartments between the plunger parts and the compartment containing the solid component ahead of the plunger parts and between the latter and the needle.\nThe bypass thus serves to allow transfer of the liquid substance from the upstream compartment to the downstream compartment for mixing with the solid substance when the plunger part separating the two compartments is shifted toward the needle-end of the barrel to an extent sufficient to allow the bypass to communicate between the compartments around this plunger part.\nThe drawback of this type of syringe, however, is that with an excessive pressure upon the stem or rod of the piston, the liquid can be transferred through the bypass with an excessive velocity so that the substance within the compartment at the downstream side of the bypass will be forced from the needle passage out of the syringe in an uncontrolled manner. This can lead to undesired contamination of the user of the syringe as well as of the patient and may be detrimental to the patient if the administration of a precise quantity of a medicament is essential."}
{"text": "1. Technical Field\nThe invention disclosed and claimed herein generally pertains to a method for applying a protective coating to at least one layer of a giant magneto-resistive (GMR) sensor, to inhibit corrosion. More particularly, the invention pertains to a method of the above type wherein a protective coating is applied to a layer of copper or copper alloy sandwiched between other layers of material, such as an alloy of cobalt and iron. Even more particularly, the invention pertains to a method of the above type wherein the copper and cobalt-iron layers are included in a stack of layers comprising a sensor for a read head of a magnetic media data storage system.\n2. Description of the Related Art\nThe continuing requirement to increase storage densities of magnetic media data storage systems, such as magnetic tape drive systems, is increasing the need for more sensitive magneto-resistive sensors. Currently, anisotropic magneto-resistive (AMR) sensors are being employed in the read heads of such systems. In AMR sensors, corrosion generally has not been a significant issue associated with their use. In these sensors corrosion has been mitigated through the use of half-power storage, which keeps the sensor material warm when not in use. This, in turn, reduces the amount of corrosive gases that are adsorbed onto the sensor material. However, the maximum change in sensor resistance (ΔR/R) for an AMR sensor is only 2%. As is known by those of skill in the art, ΔR/R is a metric of total signal available from the sensor, and thus indicates the sensitivity of the sensor.\nIn the effort to increase data storage density in magnetic media, it has been recognized that the ΔR/R of a read sensor may be significantly increased by using a device known as a giant magneto-resistive (GMR) sensor. The ΔR/R for a GMR sensor is 10-20%. However, a major drawback associated with this type of sensor has been its increased susceptibility to corrosion. A GMR sensor generally comprises a stack of layers, wherein a central layer is formed of copper or a copper alloy. Hereinafter, for convenience, the term “Cu” is used to mean or refer to the copper or copper alloy material forming such central layer. Other layers are formed of materials such as alloys of platinum-manganese (PtMn), cobalt-iron (CoFe), nickel-iron (NiFe), ruthenium (Ru) and tantalum (Ta). The most susceptible layer to corrosion is the crucial Cu layer. The Cu layer carries substantial electric current, and the GMR effect occurs at the interface of the copper and adjoining layers, which usually are comprised of a cobalt-ferrite alloy (CoFe).\nSeveral solutions to the Cu corrosion problem are available in the prior art. One such solution is to place a protective layer on top of the stack of sensor layers, to prevent corrosion of any of the layers. However, at least two problems exist with this solution: the protective layer can be worn off by the moving magnetic tape or other media, and the protective layer, when present, introduces spacing loss which can reduce the sensor signal and resolution to unacceptable levels.\nA second prior art solution is to bury the GMR sensor within the read head of the tape drive system. This is commonly done in either a yoke structure or in a flux guide structure. However, a number of problems also arise with these solutions. Signal strength is reduced, since yoke and flux guide designs have a maximum efficiency of only 50%, with efficiency usually being closer to 20-30%. Also, the manufacture of the yoke structure tends to involve substantial complexity.\nThe above prior art solutions to corrosion of the Cu layer in a GMR sensor are further described hereinafter, in connection with FIGS. 3-5. It would clearly be of great benefit to provide a process for protecting the vulnerable Cu layer of a GMR sensor that did not introduce additional spacing losses between the media and the active layer of the sensor."}
{"text": "1. Field of the Invention\nThe present invention relates to an assemble processing system and, more particularly, to an assemble processing system adapted to processing macro features.\n2. Description of the Related Art\nIn a prior art assemble processing system, a macro expansion processing requires a considerably long time to generate a macro expansion area. This will be explained later in more detail. As a result, the time for executing a macro reference processing is remarkably increased. In addition, the memory area required is also increased."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to telecommunications systems and, more particularly, to a method and apparatus for scheduling an appointment in a communications network, such as a packet network, e.g., a voice over internet protocol (VoIP) network.\n2. Description of the Related Art\nGenerally, telecommunications systems provide the ability for two or more people or machines (e.g., computerized or other electronic devices) to communicate with each other. A telecommunications system may include various networks for facilitating communication that may be generally organized into packet networks and circuit-switched networks. An exemplary circuit-switched network includes a plain old telephone system (POTS), such as the publicly switched telephone network (PSTN). Exemplary packet networks include internet protocol (IP) networks, asynchronous transfer mode (ATM) networks, frame-relay networks, and the like. One type of packet network is a voice-over-internet protocol (VoIP) network.\nSmall to mid-sized business owners are increasingly employing packet networks, e.g., VoIP networks, to utilize alternative business class telephony applications, such as IP Centrex. At present, however, enterprise customers are not able to designate call treatment preferences that can be configured into service logic that processes incoming calls for handling appointment scheduling.\nThus, there is a need in the art for a method and apparatus that enables and enterprise owner to define incoming call service logic for handling appointment scheduling."}
{"text": "To give a driver good drive feel, an engine is typically sized relative to vehicle weight using design parameters of other components, such as, for example, torque converters, transmission gear ratios, and manifold volume. In particular, manifold volume plays an important role in drive feel since when a driver requests and increase in wheel torque, conventional stoichiometric engine controls cannot instantly provide this wheel torque since these conventional engine control systems control engine torque primarily through control of a throttle. Thus, the engine must be somewhat oversized so that the driver has acceptable vehicle launch performance despite delays due to manifold volume, where vehicle launch refers primarily to operating pedal engagement at low vehicle speeds.\nThe inventors herein have recognized a disadvantage with the above approach. In particular, when an engine is somewhat oversized, fuel economy decreases due to excess weight of the somewhat oversized engine. Prior approaches for increasing an engine power to weight ratio, such as supercharges, still suffer from disadvantages due to torque increase delays, and therefore cannot improve initial vehicle launch feel."}
{"text": "Wireless IBOC signals may include information associated with one or more services, such as audio, traffic, or data services for HD Radio™. Less than all of the services in an IBOC signal, however, may be of interest to a receiver device that receives the signal. Existing receiver devices may nevertheless store all received services in a memory and subsequently determine which of the stored services are applicable to the receiver device. Such determinations may require more memory than is necessary for processing services of interest, which may be inefficient and costly. Furthermore, existing standards that prescribe formats for communicating IBOC information may not provide information that facilitates more efficient processing of IBOC information."}
{"text": "(1) Field of the Invention\nThis invention pertains to an air valve of a swimming mask which utilizes the mechanism of springs to regulate the opening and closing of the tube for air circulation on the surface of the water. The presence of the springs will prevent the tube from closing permanently to stop the air inflows. The drainage system of this swimming mask employs a flapper which can drain out any water that comes in from the upper end of the tube.\n(2) Description of the Related Art\nRepresentative of the related art known is the U.S. Pat. No. 2,317,236 to C. H. Wilsen et al in which the air intake and outlet valve is a ball valve which is to open or close the flange to prevent water from entering the tube under the surface of the water. In this drainage system, water which may enter the tube will pass through the face to the bottom valve and eventually be purged out through the bottom valve by either submerging the head underneath the water or blowing into the mouthpiece while pressing against the ball valve by hand.\nThe applicant's swimming mask is imagined according to the human respiratory system around the throat. The biological principle of separating water and air at the pharynx into esophagus and trachea is applied to the design of the protective cover and the air intake and outlet valve on the top part of the plastic tube that is connected with the mask. During swimming, the valve utilizes the springs to open or close the mobile cover depending on the absense or the presence of an external water pressure from the top of the protective cover. This valve will discriminately allow only air to enter into the mask through the plastic tube, thus giving the user a bountiful supply of air. If water enters into the protective cover accidentally, it will be drained out automatically through a hole at the base of the protective cover, preventing it from entering the plastic tube."}
{"text": "The present invention relates generally to a multiple processing unit system in which the processing units can communicate directly with each other, and more particularly to a novel communication bus and bus protocol which facilitates communication between processing units.\nMulti-processor or multi-processing unit systems are well known in the art. To fully utilize the benefits of having multiple processing units working in parallel, the processing units should communicate with one another. However, most bus structures connecting multiple processing units are designed to allow the processing units to share memory, not communicate directly. That is, the busses are memory busses, not communication busses. Memory busses are configured to handle bursts of data efficiently, such as cache lines or other data blocks. In designs in which the memory bus is used to connect processing units to DRAM memory, it also is important to perform longer bursts of data transfer for efficiency, as DRAM has a set-up overhead when accessing an address which is not adjacent to a previously accessed address. This burst nature gives the bus higher bandwidth at a cost of higher latencies.\nIn certain situations, it is important to have low latency transfers. For example, in a multi-media system performing a 3D pixel rendering function, one processing unit may be configured to process a list of geometry coordinates, while another processing unit may be configured to render the scene on a video display. With the current systems, the first processing unit writes the geometry coordinates to memory, and then the second processing unit pulls the coordinates out of memory to render the scene. With this configuration, the memory buffer decouples the two processes, creating latencies caused by writing data to and reading data from the memory. These latencies reduce the true parallel processing benefits from the multi-processing unit system, because the second processing unit is delayed in rendering the 3D scene. To more efficiently process the data, it is preferable for the first processing unit to pass the data directly to the second processing unit to avoid the memory latencies. That way, the second processing unit can perform the rendering process as soon as possible.\nThus, what is needed is a communication bus and associated protocol that facilitates low latency transfers between multiple processing units.\nIn accordance with the invention, a multi-processing unit system including a plurality of processing units in direct communication via a communication bus is presented. The system includes a communication bus arbiter having a communication packet multiplexer. Each of the processing units includes a communication bus interface comprising a transmitter interface and a receiver interface. Each of the transmitter interfaces is connected to the communication packet multiplexer of the communication bus arbiter via separate 32-bit interfaces. Each of the receiver interfaces is connected to the communication packet multiplexer of the communication bus arbiter via a single 32-bit bus.\nThe system may further comprise a first control signal connection means for communicating control signals between the transmitter interface and the communication bus arbiter and a second control signal connection means for communicating control signal between the receiver interface and the communication bus arbiter. In accordance with one embodiment of the present invention, the first control signal connection means comprises two 1-bit communication connections, one for communicating bus request signals from the transmitter interface to the communication bus arbiter, and the other for communicating bus request acknowledgment signals from the communication bus arbiter to the transmitter interface. Similarly, the second control signal connection means preferably comprises two 1-bit communication connections, one form communicating receiver buffer full signals from the receiver interface to the communication bus arbiter, and the other for communicating receive packet enable signals from the communication bus arbiter to the receiver interface.\nIn accordance with one embodiment of the present invention, the processing units preferably communicate with one another using communication packets comprising a 32-bit header and a 128-bit data packet. The 128-bit data packet preferably is presented as 4 32-bit data packets.\nThe present invention further comprises a novel communication protocol for communicating data packets between the processing units over the communication bus. In accordance with one embodiment of the protocol or the present invention, a first processing unit sends a bus request to a bus arbiter. If the communication bus is open for communication, the bus arbiter grants the first processing unit access to the bus. The first processing unit then sends a header packet to the bus arbiter, which includes a target ID. The bus arbiter takes the target ID and determines if a second processing unit associated with the target ID is able to receive data from the first processing unit. If the second processing unit is able to receive data, the first processing unit sends a 128-bit data packet to the bus arbiter, which in turn, sends the data packet to the second processing unit. In accordance with one embodiment of the present invention, the 128-bit data packet is sent in 4 separate 32-bit packets. It typically takes 5 clock cycles to send the 128 bits of data to the second processing unit."}
{"text": "The voice message service means that a user could record his voice with a mobile terminal such as a cell phone and send the recorded voice to one or more friends for listening. Furthermore, the user may also perform operations, such as reception, forwarding, inquiry, reply or playing on-demand of a voice message, according to the prompt tone of the mobile terminal, for example a cell phone. Voice messages compensates for the difficulty in delivering voice and inconvenience in information input for traditional text messages, therefore solving the problem that those persons who do not use the short messaging service because of being unfamiliar with Chinese phonetic alphabet cannot communicate with others with messages.\nWhereas, since a voice message received by a user is an audio file sent from a sender, it is impossible for the user to view the voice message intuitively. When a mobile terminal such as a cell phone has a large number of messages stored locally, it is extremely hard for the user to find out a specific voice message. Therefore, voice message searching is very inconvenient, which greatly reduces user experience."}
{"text": "1. Field of the Invention\nThis invention relates to an optical card handling system used with a general-purpose operating system and a method of accessing an optical card.\n2. Description of the Related Art\nInformation processing techniques have recently been making remarkable progress, and means for recording an increasingly large volume of information have been required. In this connection, what is now attracting attention is an optical information recording/reproducing apparatus.\nThis type of apparatus uses a recording medium which information can be optically written into and read from. As one of optical information recording/reproducing apparatuses of this type, an optical card apparatus using an optical card as a recording medium has been put to practical use. The optical card apparatus is capable of recording information on an optical card and reproducing the recorded information.\nThe optical card contains an information recording layer which has the property of changing irreversibly in the presence of heat. By projecting the laser light gathered by a lens onto the information recording layer, pits (holes) can be made in the layer. In this way, the presence and absence of pits (holes) made in the information recording layer represent a data string. Thus, by controlling the formation of pits according to the contents of data to be recorded, the desired data can be written. The reproduction of data is effected by optically sensing the presence and absence of pits in the information recording layer and converting into data.\nThe optical card has a recording capacity thousands to tens of thousands of times that of a conventional magnetic card. Once data is recorded in an area, the data cannot be rewritten. However, the data capacity that can be recorded on a single optical card is as large as 2 Mbyte to 4 Mbyte. Therefore, the optical card is considered to applied to a variety of uses including bankbooks, pocket maps, and prepaid cards used in shopping. Furthermore, by making use of the advantage that data cannot be rewritten, the optical card is considered to be applied to the case where data must not be falsified.\nAn example of the optical card is shown in FIG. 12. The optical card 1 shown in the figure is of the credit card size and has a data recording section 2 on its one side. The data recording section 2 is divided into a lot of tracks 2a as the magnetic recording medium is. Each track 2a is straight. Data is recorded sequentially track by track in the direction in which the track extends. The specific areas at both ends of each track 2a, that is, the specific regions at the begin and terminate ends of each track 2a, are used as ID sections 3a, 3b in which information items such as track addresses are recorded.\nAs described above, the optical card 1 is a medium on which data is recorded in the form of the presence and absence of pits, more specifically, on which data is recorded by forming pits in each track. Once pits are formed, they cannot be filled in, so that new data cannot be written over the track in which data is already written. This type of optical card is called a write-once, read-many type.\nHowever, data can be written in an unused track. Because each track has the ID sections 3a, 3b, a file can be updated as follows. When the data is updated, the whole file including the data items to be updated is recorded on an unused track, and the old data file is discarded. As with the magnetic disk, some tracks are used for file management. Once the file is updated, new file management information is written on unused ones among the file management tracks, and management is effected so that the updated file can be distinguished from the discarded file.\nTo record and reproduce data onto and from the optical card 1, a read/write beam of laser light must be swept track by track. The light beam may be swept along the track. In addition, the same result may be produced by placing the light beam at a constant position and causing the optical card to move repeatedly along the track.\nIn general, the optical card is moved along the track and the light source of the optical beam is moved across the track so that data can be read from and written into the desired track.\nA detailed structure of the data recording section 2 will be described. As shown in FIG. 13, a plurality of track guides 2a1 are provided at regular intervals in parallel with the track so that each track 2a2 may be formed between two adjacent track guides. Each track guide 2a1 can be sensed optically. A sensor for sensing the track guides is provided. Each track guide is sensed by the sensor. Because each track guide 2a1 and each track 2a2 are arranged along the same lane, the position of the light beam can be controlled so as not to depart from the track lane by sensing the track guide corresponding to the desired track with the sensor and controlling the position of the data read/write light beam so that the beam may move along the track guide. Of course, the data read/write light beam is positioned in the center of the track lane.\nIn this way, data pits 2a3 representing the recording information can be formed and read along the track 2a2. The track 2a2 of the data recording section 2 (not shown) is further divided into sectors. Data is read and written in sectors.\nThe data recording section 2 which data is read from and written into has a data area 2b1, a directory storage area 2b2, a change area 2b3, and a change information area 2b4. Namely, some of the tracks in the data recording section 2 are allocated to the data area 2b1, directory storage area 2b2, change area 2b3, and change information area 2b4, which are used as dedicated areas.\nThe change area 2b3 is an area used for changing the data when a recording error has occurred due to a defective track on the optical card 1, or for logically rewriting the data. The change information area 2b4 is an area in which the fact that the data is written in the change area 2b3 for replacement is recorded.\nThe data area 2b1 is an area in which the contents of the data forming a file is recorded. The directory storage area 2b2 is an area in which file management information is stored. In the directory storage area 2b2, the necessary information for file management, including a file name, the locations in which the file is recorded, and the length of the file, is recorded. Therefore, by searching the directory storage area 2b2 for the data storage location, the contents of the data can be retrieved.\nWith the optical card, the reading and writing of information is effected by an optical head composed of an optical system as shown in FIG. 15.\nIn the optical head shown in FIG. 15, the laser light emitted from a light-emitting element 4 is shaped by a collimate lens into parallel light and then passes through an object lens 6. The object lens 6 focuses the light onto the optical card 1. The reflected light from the optical card 1 passes through a mirror 7, an image forming lens 8, and a beam splitter 9, and is supplied to a detector 10. The detector 10 senses the reflected light and converts it into an electric light.\nThe detector 10 produces an electric signal according to the presence and absence of pits in the track. Thus, moving the optical card 1 along the track lane enables the data pits 2a3 shown in FIG. 13 to be sensed by the detector 10 in the form of the strength and weakness of the electric signal. Furthermore, in the optical head, the light obtained from the beam splitter 9 is supplied to a focus sensor 11 to sense the deviation of focus. By actuating an object lens driving section 12 according to the sense result of the focus deviation, the position of the object lens 6 is adjusted and thereby the focus is adjusted. To reciprocate the optical card 1 in its longitudinal direction (along the track lane), an optical card driving section 13 is provided. Explanation of a tracking error signal sensing system will be omitted.\nWhen the optical card 1 is used in a computer system, attention should be given to the following points. In the computer system, to use the optical card 1, an optical card recording/reproducing apparatus is connected via an interface. The optical card recording/reproducing apparatus is controlled in various ways under the control of the operating system (OS), the basic program for the computer system, and reads and writes data from and onto the optical card 1.\nMany personal computer systems and workstations use UNIX (a trademark of AT&T), MS-DOS (a trademark of Microsoft Corporation), Windows (a trademark of Microsoft Corporation), Macintosh-OS (a trademark of Apple Computer, Inc.), and DOS/V (a trademark of IBM corporation) as an operating system (OS). Under the present conditions, these OSes occupy the position of standard OSes in the fields of personal computer systems and workstations.\nWith these standard OSes, to manage files for an external storage apparatus, various pieces of information on files are stored in the directory, and on the basis of the information in the directory, the reading and writing of files is effected. These standard OSes differ from each other in the format of directory information and are not interchangeable with one another. Therefore, in a computer system using these standard OSes, when the optical card is read and written, access cannot be effected unless the data format in which data is recorded in the directory storage area 2b2 on the optical card coincides with the data format of the directory on the OS side.\nThat is, once the format of the directory storage area 2b2 on the optical card 1 is determined, operating systems which can access the optical card 1 are limited.\nGenerally, the optical card 11 is often carried with the user and used in various places. Consequently, he may encounter the case where his card must be read from or written into with an optical card recording/reproducing apparatus of a computer system running on an OS different from the OS used in recording the data on the optical card."}
{"text": "This invention relates to a process for preparing trifluoromethyl-nitrodiphenyl ethers having the formula ##STR4## wherein m in 1 to 3, n is 1 or 2 and R is hydrogen, carboxyl, carboxylate, formyl, keto, oxime, cyano, alkyl, alkoxy, chloro, bromo, or N, N-dialkylamino. The trifluoromethylphenyl-nitrophenylethers that may be prepared by the process of this invention are members of a class of compounds useful as herbicides and/or as intermediates for the preparation of various herbicides, pesticides, dyestuffs, and pharmaceuticals. In recent years, the development of commercial utility for trifluoromethylphenyl-nitrophenylethers in the agricultural and pharmaceutical fields has led to considerable activity in the investigation of methods of preparation. The following U.S. patents disclose the preparation and/or use of various trifluoromethylphenyl-nitrophenylethers: U.S. Pat. Nos. 4,262,152; 4,031,131; 4,259,510; 3,941,830; 3,798,276; 3,928,416; 3,784,635; 4,063,929; 4,087,272; 4,263,277; 4,263,041; 4,002,662; 4,001,005; and 3,979,437.\nVarious methods for the preparation of 2-chloro-4-trifluoromethyl-4.sup.1 -nitrodiphenyl ethers are known. One method, disclosed in U.S. Pat. No. 4,259,510 comprises reacting a 2-chloro-4-trifluoromethylphenolate with a p-nitro-halobenzene. However, the phenolate reactant is not readily available commercially and, as a result, this synthesis route involves a preliminary step to prepare the phenolate. In another method, suggested in U.S. Pat. No. 4,031,131, a 3-chloro-4-halobenzotrifluoride is reacted with a phenoxide and the resultant diphenyl ether is subsequently nitrated. In addition to requiring an additional step, the subsequent nitration is inefficient and results in the production of undesired isomers. To avoid the disadvantage of a subsequent nitration step, the 3-chloro-4-halobenzotrifluoride may be reacted with a p-nitrophenoxide. However, it has been found that the presence of the nitro group drastically reduces the nucleophilicity of the phenoxide oxygen that is the reactive site, so that reaction with compounds such as 3,4-dichlorobenzotrifluoride are found to be extremely inefficient.\nNevertheless, it has now been found that, despite the adverse effect of the nitro group on the nucleophilicity of the phenoxide oxygen when the chlorohalobenzotrifluoride reactant is a chloro-fluorobenzotrifluoride such as 3-chloro-4-fluorobenzotrifluoride, the reaction with nitrophenoxides, especially p-nitrophenoxides can be easily affected, since the reaction of the phenoxide oxygen at the fluorine site proceeds smoothly even without a catalyst to provide the trifluoromethylphenyl-nitrophenylether product in high yield. In this matter, the need for a subsequent nitration step is eliminated. Furthermore, the method of this invention, utilizing chloro-fluorobenzotrifluoride reactants is particularly suitable for reactions with nitrophenoxides carrying highly sensitive functional groups, such as aldehydes which might otherwise be further oxidized in a subsequent nitration step such as required by the prior art process described above.\nChloro-fluorobenzotrifluorides may be prepared by the vapor phase chloro-denitration reaction of the corresponding chloro-nitrobenzotrifluoride with a chlorinating agent. Thus, for example, 3-chloro-4-fluorobenzotrifluoride may be prepared by vapor phase chloro-denitration reaction of chlorine with 4-fluoro-3-nitrobenzotrifluoride. The chloro-denitration process is carried out under conditions of temperature and pressure appropriate for a vapor phase reaction, the exact conditions being dependent on the properties of the particular reactants employed. Typically, the process is carried out at atmospheric conditions and at a temperature in the range of about 250.degree. to about 450.degree. Celsius, preferably 300.degree. to 400.degree. C. Typical of the chloro-substituted benzotrifluorides that may be prepared in this manner and employed as reactants in the process of the present invention are 3-chloro-4-fluorobenzotrifluoride; 4-chloro-3-fluorobenzotrifluoride; 2-chloro-5-fluorobenzotrifluoride; 5-chloro-2-fluorobenzotrifluoride; 2-chloro-4-fluorobenzotrifluoride; 4-chloro-2-fluorobenzotrifluoride; 3-chloro-5-fluorobenzotrifluoride; 2,5-dichloro-4-fluorobenzotrifluoride; 4,5-dichloro-2-fluorobenzotrifluoride; 3,5-dichloro-4-fluorobenzotrifluoride; 3,4-difluoro-5-chlorobenzotrifluoride; 2,5-difluoro-3-chlorobenzotrifluoride; 3,5-difluoro-4-chlorobenzotrifluoride; and the like."}
{"text": "The present invention pertains to internal mold release compositions and to polyol compositions containing same and process for preparing polymeric products.\nPolyether polyurethane moldings are being increasingly used in the manufacture of automobiles, furniture and in home construction. Molded polyether polyurethanes are especially important because they are lightweight and are resistant to moisture, weather, temperature extremes, and aging. As an illustration, molded polyether polyurethane elastomers have become of special interest in the manufacture of force-reducing impact media such as safety impact bumpers for automotive vehicles and impact resistant automotive fascia.\nThe high demand for molded polyether polyurethane requires articles requires that they be produced in the largest numbers in the shortest possible time. Polyurethane-forming mixtures are eminently suited for mass production because the reactants are liquid, that is they are pumpable, and are quick-reacting. The problem has existed, however, in providing adequate mold release in the shortest possible time to take fullest advantage of the unique capabilities of the polyurethane systems.\nHeretofore, release of molded articles from molds in which they have been formed has been achieved by coating the surface of the mold cavity with an agent which facilitates release of the molded article from the walls of the mold cavity. Procedures such as this are described in U.S. Pat. Nos. 3,694,530, 3,640,769, 3,624,190, 3,607,397 and 3,413,390. This method has certain disadvantages. The agent, after molding, adheres to the surface of the molded article thereby removing such from the surface of the mold. As the mold release agent is removed from the mold surface, it must therefore be replaced so as to provide continued release of the molded articles from the mold. The necessity for repeated additions of mold release agent results in higher costs due to low productivity as a result of the additional time incurred in applying such additional quantities of mold release agents to the mold surfaces.\nIn addition, mold build-up may become a problem, since a fine film of urethane is left in spot areas of the mold surface. This build-up on the surface of the mold cavity walls eventually covers and obscures any detail on the mold cavity surface desired to be imparted to the molded article. Also, the presence of the release agent adhering to the surface of the molded article can impede subsequent operations on the article, such as painting or adhering operations.\nAdditionally, the need to reapply the release agent after each molding or a limited number of moldings interrupts the molding operation and slows down output.\nThe use of internal mold release agents for use in molding polyurethane articles has been disclosed by Boden et al in U.S. Pat. No. 3,726,952, Godlewski in U.S. Pat. No. 4,024,088, Bonin et al in U.S. Pat. No, 4,098,731, Sparrow et al in U.S. Pat. No. 4,130,698, Godlewski in U.S. Pat. No. 4,111,861, Kleimann et al in U.S. Pat. No. 4,201,847 and Godlewski in U.S. Pat. No. 4,220,727.\nSome of these internal mold release agents bleed or creep to the surface of the molded article. Some of these articles can not be painted even after appropriate preparation steps for painting has been done. Others are incompatible with polyether polyols. Most of them seriously reduce the activity of the catalyst. Almost all show degradation of physical properties such as reduced elongation.\nThe use of the \"salts\" described in U.S. Pat. No. 3,726,952 have not been effective release agents for reaction injection molding (RIM). While showing release characteristics per se their use has demonstrated in a screening program wherein hand mixed formulations are cast into an open mold other serious problems, namely: (1) degradation of the tin catalyst employed in the formulation, (2) excessively long gel and cure time, and (3) poor physical properties. These problems are believed to be caused by the presence of free carboxylic acid. It is released from the salt by the reaction of the amine with the isocyanate, and it is believed that the presence of these free carboxylic acids, or any acid, interferes with the cure rate of the hydroxyl-isocyanate reaction to form a urethane structure as disclosed in J. Polymer Science, Polymer Chemistry Edition, Vol. 19, 381-388 (1981) John Wiley & Son, Inc.\nThe reactivity or catalyst kill problem can be overcome to a certain degree by using tertiary amines in place of primary or secondary amines. Both U.S. Pat. Nos. 3,726,952 and 4,098,731 describe this technique. Since isocyanates cannot react with tertiary amines the salt cannot be split; it thus remains neutral (the carboxylic acid is not free), hence, catalyst kill does not seem evident. The use of tertiary amines, however, often shows bleed out or exudation problems which in turn result in poor paint adhesion. Further, retention of physical properties is seldom possible because of either excess reactivity when using very catalytically active amines, or because of plasticizer effects brought about by excessively long tertiary amine molecules.\nThe technology of U.S. Pat. No. 4,111,861 states that polar metal compounds can be employed to overcome catalyst kill problems brought on by the presence of fatty carboxylic acids. It states that metal ions must be present in an amount sufficient to neutralize the acid. Reference is made to the use of the Bi, Pb, Na, Li, and K ion, with sodium carbonate, sodium oleate, and potassium laurate being exemplified. They also show sodium oleate alone to be an effective release agent. When evaluated in RIM polyol systems as a single additive, it failed to show adequate release characteristics in a screening program wherein hand mixed formulations were cast into an open mold.\nZinc stearate has long been known to be an effective release agent for most thermoplastics. It is also used in polyester sheet molding compounds. When evaluated in RIM polyol systems containing only hydroxyl groups as the active hydrogen-containing source, zinc stearate as a single additive failed to show adequate release characteristics in a screening program wherein hand mixed formulations were cast into an open mold. Zinc stearate was observed to dissolve in a mixture of oleoyl sarcosine and excess polyoxypropylene diamine of 400 MW and the resultant mixture performed as an effective mold release agent.\nThe present invention provides for an improvement in one or more of the following: (1) increased multiple release, (2) increased ease of release, (3) effective and very stable catalyst reactivity, and (4) minimally altered physical properties in molded parts."}
{"text": "This invention relates to optical detectors useful for lightwave communications, and in particular to a monolithic silicon photodetector which includes means for aligning an optical fiber for edge illumination of the detector.\nLightwave communication systems are currently enjoying an intensive development effort. In the present technology generally, silicon photodetectors are utilized to detect light produced by GaAlAs lasers at a wavelength of approximately 0.85 microns. The photodetectors are usually p-i-n or avalanche diodes which are illuminated through the top surface of the device (see, for example, Melchior, \"Detectors for Lightwave Communication,\" Physics Today, Volume 30, pages 32-39 (November 1977)). While such systems are more than adequate, it is recognized that an increase in distance between repeaters over that presently available is highly desirable for more economical systems. This presently requires use of light at longer wavelengths.\nFiber attenuation in available fibers decreases as the wavelength of light increases and it is therefore desirable to provide devices which can detect the longer wavelengths. Presently available top-illuminated photodetectors are generally limited to the 0.8-0.9.mu. range since the detectable light absorbing region, which is defined by the depth of the depletion region in the device, is shallow and does not provide sufficient distance for light of longer wavelengths to be absorbed therein. That is, at wavelengths beyond 0.9 .mu.m, the quantum efficiency of conventional top-illuminated silicon detectors drops off rapidly due to the decrease in absorption coefficient with wavelength (see, e.g., Conradi, \"Fiber Optical Transmission Between 0.8 and 1.4 .mu.m\", IEEE Transactions on Electron Devices, Vol. ED-25, No. 2, pp. 180-191 (February 1978)).\nIf one wishes to provide detectors which can operate above 0.9 .mu.m, one has basically two alternatives: to construct special silicon detectors capable of absorbing and detecting light at longer wavelengths or to develop detectors in other material such as germanium and alloys of III-V compounds. Considerable effort is now being made in the latter approach since fiber attenuation in available fibers is minimized at a wavelength of approximately 1.3 .mu.m and absorption in silicon is limited to approximately 1.1 .mu.m due to its small bandgap (1.12 eV).\nThere are, however, definite advantages in retaining silicon as a photodetector material. These include the fact that silicon technology is well developed, the material has a greater inherent sensitivity than other materials being explored, and it permits integrating the detector with other elements on the same chip. It is possible to construct silicon photodetectors which are responsive to wavelengths of 1.0-1.1 .mu.m and therefore permit lower fiber attenuation while retaining silicon as the detector material. This construction basically involves making the optical path length longer than the diode depletion depth to accommodate the absorption requirements of the longer wavelengths. This can be accomplished, for example, by illuminating the detector at an edge so that light travels parallel to the device surface. This defines the detectable light absorbing region by the length of the depletion region which is considerably greater than the depth. However, since the thickness of the depletion region in the detector is very small, the vertical alignment of the fiber becomes critical. Alignment by normal mechanical means therefore becomes very difficult.\nIt is therefore a primary object of the invention to provide a silicon detector structure which is responsive to light at long wavelengths and includes means for precise alignment of the fiber with the detectable light absorbing region of the detector."}
{"text": "The present invention relates to a discharge lamp apparatus that uses a discharge lamp as a light source, and particularly to an apparatus in which an electronic controller unit for applying a voltage to the discharge lamp is directly coupled with the discharge lamp.\nA discharge lamp apparatus that uses a discharge lamp as a light source is used as a vehicle headlight. An electronic controller unit that generates and controls a voltage applied to the discharge lamp includes a DC/DC converter for transforming an output voltage by switching an input voltage by a power device, a high voltage generation circuit for generating, from the output voltage of the DC/DC converter, a high voltage applied when lighting of the discharge lamp is initiated, and the like.\nWhen the temperature of electronic circuit components forming the controller unit rises, it is likely that, for instance, the soldered part of the circuit components melt and the circuit components operate erroneously. In the discharge lamp apparatus disclosed in JP-A-2000-235809, a part of a metallic heat radiator thermally coupled with a circuit substrate mounting circuit components thereon is exposed outside a headlight so that the heat generated by the discharge lamp and the circuit components may be radiated to the outside of the headlight through the metallic heat radiator.\nIn the discharge lamp apparatus disclosed in JP-A-2000-235809, an igniter part and a lighting device need be connected by a harness. This increases the number of component parts, complicates assembling work and adds manufacturing cost.\nIt is therefore proposed to directly couple and electrically connect the discharge lamp and the controller unit. However, if the discharge lamp and the controller unit are directly coupled and the controller unit is disposed near the discharge lamp, the internal temperature of the controller unit rises due to heat transferred or radiated from the discharge lamp and the heat generated by the controller unit itself. This is likely to cause erroneous operation of the circuit components in the controller unit.\nIt is therefore an object of the present invention to provide a discharge lamp apparatus that uses no high voltage wire nor high voltage connector and suppress rise of temperature of an electronic controller unit.\nIn a discharge lamp apparatus according to the present invention, a discharge lamp and an electronic controller unit for applying a voltage to the discharge lamp are directly coupled with and electrically connected to the discharge lamp. Therefore, a high voltage connector and a high voltage wire for connecting the discharge lamp and the controller unit are not necessitated.\nPreferably, a second casing mounting a power device of a DC/DC converter has a thermal conductivity higher than that of a first casing coupled with the discharge lamp. As a result, heat generated by the discharge lamp is less likely to be transferred from the first casing to the second casing, and the heat of the discharge lamp is less likely to be transferred circuit components mounted in the second casing.\nIn addition, heat generated by the power device of the DC/DC converter is more likely to be radiated to the outside of the second casing from the second casing that has the thermal conductivity higher than that of the first casing. Because the power device of the DC/DC converter generate more heat among the controller unit, the heat generated by the power device is readily radiated from the second casing to the outside of the second casing, thus suppressing rise of temperature of the circuit components in the controller unit including the power device. Thus, erroneous operation of the circuit components is prevented."}
{"text": "The present invention relates to the field of arithmetic processors and more particularly to a method and circuit for multiplying large numbers using a booth encoder with iterative addition.\nAs is well known multiplication of two numbers can be performed by a series of repeated additions, where the number to be added is the multiplicand and the number of times that it is added is the multiplier, the result is the product. Each step in the series of repeated additions generates a partial product. As it may be well appreciated, this process may be extremely slow when performed in a general-purpose processor, taking many clock cycles. In general terms multiplication of an M-bit number by an N-bit number will result in a product which is M+N bits long.\nThe increased need for high speed processing of large numbers has been precipitated by, for example, various cryptographic applications that use large numbers as encryption keys. These keys are at least 1024-bits long. Accordingly, multiplication using the repeated addition method that is suggested by the arithmetic definition above is often replaced by more efficient algorithms that make use of positional number representation. In general, the execution speed of an arithmetic operation is a function of two factors. One is a circuit technology and the other is the algorithm used.\nThe multiplication operation may be thought of as having two parts. The first part is dedicated to the generation of partial products and the second one collects and sums the partial products to obtain the final result. The booth algorithm is often used in the generation part because it reduces the number of partial products. The collection of the partial products can then be made using a regular array, a Wallace tree or a binary tree.\nFor an N-bit multiplier, no more than N/2 partial products are created. However, when the partial products from one multiplier operation are first added, the result is initially in a redundant format, such as carry-sum. That it, the result takes the form of two rows of binary information, a carry row and a sum row. But before this result may be used again, it must first be processed by an adder to put it into binary format. In other words, the carry row and the sum row must be added first. As mentioned earlier a Wallace tree compression unit also known in the art may be used to take the partial products and using rows of carry-save adders compress the partial products into two rows, a sum row and a carry row. A conventional N-bit full adder may then be used to add the sum row and the carry row.\nHowever this full adder is slow since a carry generated in the low order bits may ripple all the way through to the high order bits. Thus the high order must wait for the carry to ripple through all N-bits. This problem is exacerbated when large width operands are used. Alternatively, carry look-ahead and carry select adders may be used to avoid large propagation delays, but are still slow. The complexity of such adder circuits is directly related to the width of the adder. A 32-bit adder is reasonable to implement using most technologies. A 64-bit adder is extremely large, extremely slow or both.\nU.S. Pat. No. 5,944,776 describes one possible approach to eliminating the need for a full adder; by using a multiplicity of interconnected logic cells that produces a Booth output that is the Booth encoded form of the sum of a sum row and a carry row. This technique is not easily adaptable to wide width operands. Furthermore the technique described by this patent is more applicable to iterative multiplication algorithms used in a multiplicative divider.\nIn U.S. Pat. No. 5,724,280 a Booth multiplier using carry look ahead adders performs a multiplication operation in three stages: First the operands are loaded from a data bus, which takes a minimum of 8 clock cycles for a 256 -bit operand and a 32 bit bus; second while loading the second operand, Booth encoding is begun 4-bits at a time and encoded values are accumulated, which takes 64 clock cycles; and third performing a final addition on 32-bit segments while outputting a result to the data bus, which requires 8 further clock cycles when using a 32-bit adder and a 32-bit data bus. Hence, a total number of 80 clock cycles are needed with a 32-bit data bus and a 32-bit adder. This circuit would thus be unacceptably slow when used with wider width operands.\nIt is possible to implement a high speed wide-width multiplier (256-bit by 256-bit) by using a number of 256-by-256 Booth multipliers which are pipelined to obtain the desired speed. But, such a circuit would be impractically large, particularly for implementation in an ASIC.\nFurthermore, wide-width numbers can be processed by segmenting the operands and processing each of the segmented operands in a multiplier. The results from the processed segments are combined to obtain a final result. A problem associated with this technique is that carries generated while processing each segment have to be properly processed in order to obtain the correct final result, thus placing a constraint on the adder circuit, used to combine the results.\nIt is thus desirable to have a multiplier circuit using Booth encoding that performs wide width number multiplication in relatively few clock cycles and which occupies minimal chip area. Furthermore it is also desirable to have a more efficient adder circuit for processing the final result from the Booth encoder and which manages the carries generated in the Booth multiplier.\nIn accordance with this invention there is described a multiplier for computing a final product of a first operand and a second operand comprising:\n(a) a multiplier array for forming a product of the first operand and second operand in carry-save form;\n(b) a carry-save adder for adding said carry-save partial products and an accumulatd sum to produce a carry and save values;\n(c) a carry-lookahead adder for adding said carry and save values to produce a product value and a carry-out value;\n(d) a general purpose adder for adding said carry-out and said product value to produce said final product."}
{"text": "The activity of cells is regulated by external signals that stimulate or inhibit intracellular events. The process by which stimulatory or inhibitory signals are transmitted into and within a cell to elicit an intracellular response is referred to as signal transduction. Proper signal transduction is essential for proper cellular function. Defects in various components of signal transduction pathways, from cell surface receptors to activators of gene transcription, account for a vast number of diseases, including numerous forms of cancer, vascular diseases and neuronal diseases.\nProtein kinases are enzymes involved in signal transduction which phosphorylate other proteins and/or themselves (autophosphorylation). Protein kinases involved in signal transduction in eukaryotic cells can be divided into three major groups based upon their substrate utilization: protein-tyrosine specific kinases (which phosphorylate substrates on tyrosine residues), protein-serine/threonine specific kinases (which phosphorylate substrates on serine and/or threonine residues) and dual-specificity kinases (which phosphorylate substrates on tyrosine, serine and/or threonine residues). Well over a hundred protein kinases have been identified to date and more are being identified at a very fast rate (about 10 to 30 new kinases per year) through genetic and molecular biological approaches. There are conserved regions among all protein kinases, suggesting evolution from a common ancestor. For reviews on protein kinases see Kemp, B. E. (ed.), (1990) Peptides and Protein Phosphorylation, CRC Press Inc. and Hanks et al. (1988) Science 241:42-52.\nIn order to insure fidelity in intracellular signal transduction cascades it is essential that each protein kinase have exquisite specificity in downstream targets. In some cases a kinase may have only a single substrate in a cell, but in general kinases appear to have a collection of targets that allow branching of an initial signal delivered to a cell in multiple directions in order to coordinate a set of events that occur in parallel for a given cellular response (see Roach, P. J. (1991) J. Biol. Chem. 266:14139-14142). The substrate specificity of a protein kinase can be influenced by at least three general mechanisms that depend on the overall structure of the enzyme: 1) specific domains in certain protein kinases target them to specific locations in the cell, thereby restricting their substrate availability; 2) other domains in the kinase, distinct from the catalytic domain, may provide high affinity association with either the substrate or an adapter molecule that presents the substrate to the kinase; and 3) specificity is ultimately provided by the structure of the catalytic site of the protein kinase.\nAlthough the number of protein kinases that have been implicated in intracellular signaling is quite large, in only a few cases is there information about the sequence specificity of these kinases. This information has usually come from locating the phosphorylation sites on in vivo and/or in vitro substrates of the kinase. (For examples see Taylor et al., (1990) Ann. Rev. Biochem. 59:971-1005; Cheng, et al. (1991) J. Biol. Chem. 266:17919-17925; Walsh et al. (1990) in Peptides and Protein Phosphorylation. (B .E. Kemp, ed.) CRC Press Inc.pp. 43-84; Gaehlen and Harrison (1990) in Peptides and Protein Phosphorylation. (B. E. Kemp, ed.) CRC Press Inc. pp. 239-254). By comparing the sequences of several phosphorylation sites for a kinase a consensus sequence for substrates of that kinase can be determined. These types of studies have demonstrated the importance of the primary amino acid sequence around the site of phosphorylation in determining the in vivo specificity of protein kinases. Synthetic peptides can be constructed based upon the consensus sequence motif of a known phosphorylation site and individual amino acids can then be replaced one by one in order to determine the importance of particular amino acids on the K.sub.M or V.sub.max of the phosphorylation reaction.\nHowever, there are severe limitations to this approach for determining the substrate specificity of a protein kinase. The procedure is quite expensive and laborious since each amino acid residue within a phosphorylation site must be altered individually to evaluate its importance. This approach does not necessarily identify all the residues critical for substrate specificity and, furthermore, an optimal substrate sequence is not likely to be determined unless each residue is changed to every other possible amino acid residue individually and then evaluated. For example, based upon an estimation of 9 to 12 amino acid residues of a substrate contacting the active site cleft of a kinase there would be approximately 1.024.times.10.sup.13 (20.sup.10) distinct peptides to consider. Moreover, in many cases this approach is not feasible because in vivo substrates for some kinase cannot be determined with certainty. For example, if an extracellular signal activates multiple kinases which phosphorylate multiple substrates or if a signal activates a cascade of kinases it is difficult to determine which kinase phosphorylates which substrate. An even more difficult problem in determining the substrate specificity of certain kinases is that their critical in vivo substrates are often proteins which are present in vivo in very low abundance and thus are not easily detectable by in vivo phosphorylation assays. Even relatively abundant substrates can be overlooked because of their high rates of dephosphorylation by closely-associated phosphatases.\nThus, there is a need for an alternative method to that of isolating and examining the sequences of native substrates for determining the substrate specificity of a protein kinase. Consensus sequence motifs for the phosphorylation sites of many known protein kinases have not yet been determined because of the limitations to current approaches discussed above. Information on the substrate specificity of each protein kinase involved in signal transduction would provide insight into signal transduction mechanisms and could allow for the design of novel therapeutic agents based on the substrate specificities of different protein kinases."}
{"text": "The present disclosure relates to a piezoelectric vibration member, a method of manufacturing the same, and a piezoelectric vibrator.\nA piezoelectric vibrator is an apparatus generating vibrations having a certain frequency through a piezoelectric vibration member vibrating due to a piezoelectric phenomenon occurring in the piezoelectric vibration member when voltage is applied thereto.\nSince the piezoelectric vibrator has a stable vibrational frequency, such a device has been used in several core components providing a reference signal level as well as in an oscillation circuit of a computer or a communications device.\nA piezoelectric vibration member formed of crystal includes a vibration substrate using the crystal as a base material and an electrode disposed on the vibration substrate, wherein the vibration substrate may have various shapes depending on required physical properties.\nWhen a piezoelectric vibration member in a thickness shear vibration mode is formed so that a thickness of the vibration substrate is gradually reduced from a central portion thereof toward end portions thereof, a damping amount of vibration displacement in the end portions is increased, whereby an effect in which vibration energy is trapped in the central portion of the piezoelectric vibration member may be improved and frequency characteristics such as a CI value, a Q value, and the like, may be improved. An example of a shape of the piezoelectric vibration member that may accomplish a vibration energy trapping effect may include a convex shape in which a convex curved surface is formed as a main surface, a bevel shape in which a space between a flat and thick central portion and an edge of an end portion is formed as an inclined surface, a mesa shape in which a central portion is flat and a surrounding portion of the central portion is thin, and the like.\nIn a case in which the thicknesses of the thick, flat central portion and the thin surrounding portion are rapidly changed in the mesh shaped piezoelectric vibration member, an energy trapping effect may be decreased, and electrode connectivity may deteriorate. Therefore, a piezoelectric vibration member capable of solving these problems has been required."}
{"text": "To stabilize concentrated surfactant solutions utilized for enhanced oil recovery that include active surfactant content ranging from 35% to 90%, alcohols such as isopropyl alcohol (IPA) or isobutyl alcohol (IBA) are added to the concentrated solutions to stabilize them. It is preferable to blend these surfactant solutions at a convenient industrial facility rather than in the field, at a location proximate to the geographic region of use. In order to avoid having to blend the surfactant solution in the field, however, suitable storage during transportation and while awaiting use is necessary.\nThere is, therefore, a need in the art for improved surfactant solution storage."}
{"text": "Lasers have many uses, including detecting the presence or absence of objects by detecting when a laser beam is received by a sensor. A laser transmitter may send the laser beam across a void to a sensor. When the laser beam is received by the sensor, there is no opaque object in between. When the laser beam is not received by the sensor, an object may be present. Another reason why the laser beam may not be received may be that the laser or its sensor have been misaligned or maladjusted.\nA fixed laser beam may sense the presence or absence of an object using a straight beam of light. Such a sensor only detects the presence or absence along a line from the laser source to the sensor. A scanning or moving laser may be deployed to scan a planar or three dimensional area, provided that a sensor is able to detect the laser beam in the areas being scanned.\nThe alignment and calibration of a laser transmitter and receiver can be a tedious operation in some cases. In some hostile environments, the lasers or their sensors may be bumped, moved, or otherwise may be misaligned, which may inadvertently cause the laser to not be received. When such a condition is found, a technician may be dispatched to realign and test the system, then place it back in service."}
{"text": "1. Technical Field\nThe present invention relates to a resonation device, and an electronic apparatus and a moving object each equipped with a resonation device.\n2. Related Art\nSince quartz crystal oscillators vary in frequency with a variation in temperature, there are used constant-temperature quartz crystal oscillators in which the temperature of the quartz crystal resonator is maintained at a constant level by heating the quartz crystal resonator using a heating element such as a heater. In such a constant-temperature crystal oscillator, if the quartz crystal resonator is directly mounted on a substrate constituting a part of a thermostatic chamber, the heat of the quartz crystal resonator and the heating element is released to the substrate. Therefore, it becomes difficult to keep the temperature of the quartz crystal resonator constant. Further, if it is attempted to keep the temperature of the quartz crystal resonator constant, a larger amount of power must be supplied to the heating element, and therefore, there is a problem that the power consumption increases.\nIn order to solve the problem described above, in JP-A-2007-6270 (Document 1), there is disclosed a structure in which a gap is disposed between the substrate and the quartz crystal resonator in connecting/fixing the quartz crystal resonator provided with the heating element and the substrate to each other with a lead wire.\nIn such a configuration according to Document 1, the gap is disposed between the substrate and the quartz crystal resonator to thereby inhibit the heat of the quartz crystal resonator from being released to the substrate. However, since the lead wire performs both of the electrical connection and the mechanical fixation between the quartz crystal resonator and the substrate, the connection between the quartz crystal resonator and the substrate is achieved by the plate-like lead wire having mechanical strength and silver solder. Since the lead wire and the silver solder are high in thermal conductivity, the heat is released from the quartz crystal resonator to the substrate via the lead wire. Therefore, since the heat is easily released from the quartz crystal resonator, the power consumption increases, and it is difficult to reduce the power consumption."}
{"text": "This invention relates to a plasma CVD method used in the manufacture of a semiconductor device and to a semiconductor device having a metal film formed by this method, and particularly to a plasma CVD method with which a metal film formed has a smooth surface morphology and the amount of residual halogen element contained in the metal is small and a highly reliable semiconductor device having a metal film formed by this method.\nAs design rules of semiconductor devices such as LSI devices are established to cover less dimensions reduced from half micron to quarter micron or lower levels and multilayer interconnection structures come to be used in large numbers, aspect ratios of connection holes for connecting interconnection layers together have also been becoming larger. For example, in a semiconductor device based on 0.18 .mu.m rule, because with respect to a connection hole opening diameter of 0.2 .mu.m the thickness of an interlayer insulating film is about 1.0 .mu.m, the aspect ratio reaches 5. To achieve a highly reliable multilayer interconnection structure using connection holes of such small size and high aspect ratio, methods are coming to be employed wherein a metal film of Ti or the like for an ohmic contact and a TiN film or the like as a barrier metal for preventing diffusion of interconnection material are formed thinly and conformally inside the connection hole and then upper layer interconnections and contact plugs are formed by high-temperature sputtering of an Al-based metal or CVD of W.\nNormally, Ti metal films and TiN films have been formed by sputtering or reactive sputtering with bulk Ti metal as the target material. Among such methods, collimated sputtering, wherein the component of perpendicularly incident sputtered particles is increased, disclosed for example in Japanese Unexamined Patent Publication No. 6-140359, and long range sputtering, wherein a long target distance is used, have been receiving particular attention. It has been confirmed that with these sputtering methods it is possible to obtain reduced contact resistance and improved barrier characteristics compared with conventional sputtering methods. However, with these sputtering methods, because the component of sputtered particles perpendicularly incident on the substrate is increased, parts where the film thickness is extremely thin are inevitably formed around the shoulders of small connection holes having large aspect ratios and around the edges of the bottoms of connection holes. In this case, when for example blanket CVD of W is carried out in a subsequent step, the process gas WF.sub.6 permeates through the thin parts of the film and problems such as erosion of the base material layer, abnormal growth of W or the Ti metal film or TiN film flaking off occur.\nTo solve problems of step coverage not solved by these sputtering methods also including collimation and the like, methods of forming conformal Ti metal films and TiN films by CVD methods utilizing chemical reactions taking place at the substrate surface are being expected.\nCVD methods for forming Ti-based material layers proposed so far can be generally classified into two types: methods using inorganic metal-halogen compounds such as TiCl.sub.4, reported for example in Proceedings of the 44th Symposium on Semiconductor & Integrated Circuit Technology p.31 (1993), and methods using organic metal compounds such as TDMAT and TDEAT, reported for example in Proc. 11th Int. IEEE VMIC, p.440 (1994).\nWith the latter methods using organic metal source gases, in many cases a Ti metal film is formed by sputtering and then a TiN film is formed by MOCVD.\nThe former type of method using an inorganic metal-halogen compound source gas has the merit that it is possible to consecutively form a Ti metal film by H.sub.2 reduction and a TiN film by adding N.sub.2 or the like to this process gas inside the same CVD reactor.\nThe reaction reducing the metal-halogen compound TiCl.sub.4 using H.sub.2 molecules is the endothermic reaction given by the following formula (1), and thermodynamically is difficult to sustain (AG is standard heat of formation). EQU TiCl.sub.4 +2H.sub.2 .fwdarw.Ti+4HCl .DELTA.G=393.3 kJ/mol (1)\nFor this reason, formation of Ti metal films by plasma CVD wherein H.sub.2 is dissociated in a plasma and reduction by H atoms and activated H species is used has been receiving attention. This reaction is the exothermic reaction shown by the following formula (2). EQU TiCl.sub.4 +4H.fwdarw.Ti+4HCl .DELTA.G-478.6 kJ/mol (2)\nWith this formation of a Ti metal film by plasma CVD, the reaction proceeds easily and a film can be formed at a relatively low temperature. In particular, plasma CVD reactors having a high density plasma source such as ECR-CVD reactors, inductively coupled plasma CVD reactors and helicon wave plasma CVD reactors are advantageous in film formation rate and uniformity.\nHowever, even when plasma CVD is carried out using one of these high-density plasma CVD reactors, depending on the film-forming conditions the Ti metal film does not grow uniformly on the substrate and granular Ti metal grows nonuniformly.\nThis problem will be described with reference to FIGS. 1A and 1B.\nThese figures are schematic sectional views showing a substrate under processing having a connection hole 3 formed in an interlayer insulating film 2 on a semiconductor substrate 1 made of silicon or the like, and illustrate a problem arising when a Ti metal film 4 is formed on this by plasma CVD. That is, when the reduction of TiCl.sub.4 by H atoms and activated H species is not carried out well, as shown in FIG. 1A the surface morphology of the Ti metal film 4 obtained is deteriorated by grains of nonuniformly grown Ti metal.\nWhen a TiN film 6 is then formed on this, because as shown in FIG. 1B the TiN film 6 grows with the same surface shape as the Ti metal film 4 below it, this film also is granular, and at parts where adjacent TiN grains make contact with each other the film cannot grow any further. In particular, a large gap 7 sometimes forms at the corner of the bottom of the connection hole 3.\nWhen in a subsequent step a layer of a refractory metal such as W is formed by CVD with WF.sub.6 or the like as a process gas on a substrate on which is formed this kind of gap 7, the WF.sub.6 passes through the gap 7 and erodes the semiconductor substrate 1 below and destroys shallow junctions (not shown). Also, when the connection hole 3 is filled in by an Al-based metal, the Al-based metal passes through the gap 7 and reacts with the semiconductor substrate 1 made of silicon or the like to form alloy spikes. In either case, serious defects such as increased leak current at the interlayer connection result.\nAlso, when the reduction of the TiCl.sub.4 by H atoms and activated H species is incomplete, precursors TiCl.sub.x (where x is an integer below 4) and chlorine atoms are taken into the Ti metal film being formed, and residual chlorine in the Ti metal film increases. As a result, an Al-based metal layer or the like formed in a later step is corroded and there is a possibility of this leading to increased contact resistance and, in extreme cases, disconnection."}
{"text": "In wireless communication networks the User Equipment (UE) attaches to the wireless network for services. During use, the UE can have to re-attach to the network for a variety of reasons, one reason is that the UE can receive a reject message from the network. When re-attachment is required, the user equipment issues an attach request to the wireless network.\nIn the 3GPP Long Term Evolution (LTE) system, attach requests are issued by the Non-Access Stratum (NAS). The NAS is the highest stratum of the control plane between the UE and the Mobility Management Entity (MME) which functions as a network control apparatus. Operation of the NAS in a LTE system is described in 3GPP TS 24.301 Version 11.4.0 (October 2012), incorporated herein by reference.\n3GPP TS 24.301 allows two modes of operation for both the UE and the MME. In a first mode a Radio Resource Control (RRC) connection which is already open can be used for subsequent NAS transactions, including attach requests. In a second mode, a new RRC connection is required for each NAS transaction. However, there is no mechanism for the UE to determine whether the MME will operate in the first mode, and re-use an existing RRC connection, or the second mode, and require a new RRC connection.\nProviding the UE makes the correct assumption about the mode used by the MME no problems occur. However, it has been found that when the UE makes an incorrect assumption about the mode used by the network a delay of around 10 seconds can be incurred in some circumstances."}
{"text": "Accurate measurements of toxic pollutants are essential to a proper assessment of air, fluid and gas quality. In particular, ambient air particulate matter less than or equal to 10 microns in size (PM10) has been recognized as matter which can be inhaled deeply into the respiratory system and which may cause adverse health effects. The greater the concentration of particulate matter in the ambient air, the greater the risk of health problems caused by the ambient air. Particulate matter measuring devices measure the mass concentration of particulate matter within ambient air, gases or fluids to determine the quality. The measuring devices use different sensing techniques to provide continuous monitoring of particulate mass concentration. Examples of sensing techniques are beta radiation attenuation and optical sensing methods. These measuring devices can provide a warning to a user when detecting a relatively low air quality based upon a relatively large particulate mass concentration within the air.\nA particulate measuring device of the beta attenuation type is used to determine a mass concentration of a desired range of sizes of particulate matter through beta attenuation. FIG. 1 depicts a particulate mass measuring device of the beta attenuation type 10 known in the art. The monitor features a nozzle 12 having an ambient airflow chamber 14 with an axial ambient airflow pathway 16 therein. The airflow chamber has a neck 18 with an ambient air inlet 20 and curved surfaces 22 extending from the neck. A beta radiation source housing 24 disposed within the ambient airflow chamber 14 houses a beta radiation source 26 used in the determination of particulate mass concentration. During operation, the ambient air 28 being sampled (or other sample) enters the chamber through the inlet and flows to a filter tape 30 beneath the beta radiation source. Typically, before the ambient air enters the ambient airflow chamber, particles greater than a select size, (e.g. particles greater than 10 microns) are separated from the rest of the sample using impaction separation. Beneath the filter tape 30 is a beta radiation detector 32. The filter tape 30 collects the particulate matter present within the ambient air over time and the remainder of the sample flows out the outlet 34. As the amount of particulate matter collected by the filter tape 30 increases, the particulate matter attenuates the beta radiation emitted from the beta radiation source 26 as detected by the detector (i.e., the beta radiation detector senses less radiation from the beta radiation source). Thus, the attenuation of the beta radiation detected by the beta detector is related to the mass of the particulate matter collected by the filter tape 30 in a specific location. A beta radiation signal produced by the beta radiation attenuation device is used to calculate a mass concentration of particulate matter within a sample over a specified period of time.\nWith reference to FIG. 2, a problem with the prior art measuring devices is that when ambient air 28 enters the ambient airflow chamber 14, particulate matter 36 from ambient air (in particular, larger diameter or massive particulate matter having a size of 10 microns or less (PM10)) may impact and be retained on the surfaces of and/or within the ambient airflow chamber, such as curved or bent surfaces 22. Particulate matter having a size of 2.5 microns or less PM2.5 may become stuck also, however, the problem of particulate loss is more pronounced in PM10 particulate matter samples. This is because larger particles having greater inertial force are not as influenced by the air sample trajectory as the smaller particles and eventually hit the curved or bent surfaces of or within the nozzle 12 and become stuck. These stuck particles 38 will then not be accounted for when mass measurements are made often resulting in an inaccurate measurement. Further, particles stuck on surfaces of the ambient airflow chamber may accumulate on the ambient airflow chamber surfaces until random, naturally occurring phenomena cause them to dislodge and generate anomalous mass measurements in the ambient air sample.\nWhat is desired is a particulate mass measuring device that provides more accurate mass measurements of particulate matter within a sample.\nWhat is desired is a method for improving the accuracy of particular mass measurements in a particulate mass measuring device."}
{"text": "The present invention relates to a bag pack and more particularly to a cartridge for holding a stack of T-shirt type plastic film bags suitable for enclosing produce and the like.\nAs is understood by household shoppers generally as well as by those skilled in the art, produce in supermarkets is normally displayed in open bins or chill cases so that each shopper can select the individual items according to his or her preferences. It is customary to provide a supply of bags or sacks for packaging the produce So that it can easily be carried to a checkout counter. The most usual arrangement for dispensing the bags is to have them wound on a roll which is then hung so that individual bags can be drawn down and torn from the roll. This is not an entirely satisfactory arrangement in most instances since the roll will often overrun instead of releasing an individual bag. The use of both hands may be required to separate an individual bag from the roll. Further, the bags dispensed are typically simple tubes open at one end which may require a wire tie or the like for satisfactory closing. Similar problems exist in bakery and delicatessen contexts.\nThe present invention provides an effective solution to these various problems and also yields a bag pack and dispenser construction which is advantageously useful in a variety of environments including convenience stores and supermarkets as well as produce, bakery and delicatessen contexts. In accordance with one aspect of the present invention, it has been recognized that T-shirt type plastic film bags constitute a much improved mechanism for packaging produce since the loop handles extending upward on either side of the mouth of the bag can be tied to retain the contents. In accordance with another aspect of the invention, there is provided a novel cartridge and dispensing system for providing easy access to individual bags in a pack.\nAmong the several objects of the present invention may be noted the provision of a novel bag pack; the provision of such a bag pack which facilitates dispensing of individual bags from the pack; the provision of such a bag pack which incorporates a cartridge for retaining a stack of T-shirt type plstic film bags; the provision of a dispensing system for plastic film bags which is easily loaded; the provision of such a dispensing system which does not create clutter; the provision of such a bag dispensing system which is of attractive appearance and which is easily operated; the provision of such a bag pack and dispensing system which are highly reliable and which are of relatively simple and inexpensive construction. Other objects and features will be in part apparent and in part pointed out hereinafter."}
{"text": "The present invention relates to an adapter for adapting a conventional cup holder, for example a cup holder in vehicles, including cars, boats, vans, trucks, or the like, so that the cup holder can hold a container with a larger or smaller diameter than the diameter of the cup holder cavity.\nConventional cup holders include a recessed cavity which has a fixed inside diameter so that the cup holder is suitable for only a small range of container or cup sizes, thus limiting the use of the cup holder. More recently, cup holders have incorporated a tiered cavity, which forms a first smaller cavity for holding smaller containers, such as soda cans or bottles, and an upper cavity which accommodates larger containers, for examples mugs or the like. However, most of these cup holders do not accommodate the larger containers, for example, super-sized containers that are available in fast food restaurants or convenient stores, such as 7-ELEVEN.TM. stores, gas stations, and the like. When traveling long distance in a vehicle, most people like the convenience of purchasing the larger volume drinks in order to extend the period between stops and, further, to economize. Often these super-sized drinks cost only a fraction more than the smaller size drinks. Furthermore, conventional cup holders provide no means to maintain the temperature of a drink placed in the cup holder. When purchasing these super-size drinks, the lack of insulation in the cup holder is even more noticeable. Super-size drinks take longer to consume and often cool or warm up well before the drink is consumed.\nConsequently, there is a need for an adapter which will enable conventional cup holders to hold larger containers and, furthermore, which may provide means to generally maintain the temperature of the drink in the container."}
{"text": "1. Field of the Invention\nThe present invention relates to a multi-component fiber optic handpiece that is particularly safe, practical and economical for use in medical laser treatment.\n2. Information Disclosure Statement\nAn optical fiber handpiece is generally used in ophthalmic, dental, and orthopedic surgical procedures where the treatment area is confined and particularly difficult to reach. Typically, laser light is transmitted from a laser source through an optical fiber to the treatment site. The optical fiber terminates proximally in a laser connector for connection to the laser source and terminates distally in a handpiece, which is manipulated by the surgeon.\nAs can be appreciated, because of the potential for spreading infection, a handpiece used during one procedure cannot be used with another patient in a subsequent procedure unless some form of sterilization is performed. Unfortunately, the handpieces that were initially introduced to the marketplace were not designed to withstand sterilization procedures and were therefore discarded after use on a single patient. However, this method is not cost-effective.\nSince that time, attention has been given to improving the design of the handpiece so that it can be resterilized. Various handpieces have been manufactured and sold which can be resterilized using gas techniques. Unfortunately, gas sterilization is a relatively time consuming and costly approach. In contrast, high temperature steam sterilization in an autoclave is faster and less expensive. However, handpieces were generally incapable of withstanding the high temperatures that are encountered in an autoclave.\nU.S. Pat. No. 5,304,172 provides a partial solution to this problem by carefully selecting materials for the handpiece body that can withstand the maximum heat that is encountered in an autoclave, which is about 132 degrees centigrade. However, the optical fibers that are incorporated into the handpiece failed or were damaged during reuse. The disclosed fiber optic probe utilizes a two-piece optic fiber to deliver laser radiation through the handpiece. The two pieces are held in place by spring-biased ferrules that force the mating ends of the fibers together. No means are disclosed to protect the optical fiber during insertion or prior to adhesively securing the optical fibers to the ferrules.\nIt has been found through research that exposure to the input face of the fiber to the cycling of the pressurized steam occurring during sterilization in an autoclave enhanced and accelerated the formation of cracks in the silica glass at the fiber surface. It has been previously reported that water molecules in the presence of cracks in glass can accelerate the breakdown of bonds. (See, xe2x80x9cThe Fracturing of Glass,xe2x80x9d Michalske and Bunker, Scientific American, December, 1987, pages 122-129). As can be appreciated, the high temperatures generated in an autoclave can force steam molecules into any microscopic cracks that are present in the fiber, accelerating the breakdown of atomic bonds. It is also believed that when the autoclave is rapidly depressurized, turbulence is created, forcing debris into the input face of the fiber thereby increasing the damage. These cracks and other imperfections lead to breakdown of the fiber during use.\nAccordingly, it would be desirable to provide a handpiece that could be sterilized in an autoclave without the complications associated with the prior art.\nOne possible solution is to remove the optical fiber prior to resterilization. However, the fiber can be damaged during insertion and an adequate means to removably lock the fiber in position would be required. Misalignment of the fiber is prevalent when the fiber is not permanently incorporated into a handpiece device, which can result in damage to the fiber during assembly and actual use of the handpiece.\nU.S. Pat. No. 5,489,205 discloses a syringe tip locking assembly comprising means for compressing a handpiece body against a removable syringe. The four part holding mechanism described may work for a metal syringe, but will not necessarily work on a glass fiber. In a holding means such as that disclosed, if the fiber were not perfectly centered prior to being locked in place, it would be possible to lock a fiber with only 3 of the 4 prongs, thereby placing unbalanced forces on the fiber. Stress crazing could result in catastrophic failure of the fiber.\nIt is therefore an object of the present invention to find a solution to the handpiece-fiber combination that is both economical and practical for medical laser treatment.\nIt is an object of the present invention to provide a handpiece-fiber combination without the complications associated with the prior art. In particular, a multi-component handpiece is described that is safe, practical and economical for medical laser treatment.\nIt is another aim of the present invention to provide a means to protect the optical fiber before and during insertion into a handpiece device.\nIt is a further aim of the present invention to describe a means to prevent longitudinal movement of an optical fiber during surgical laser therapy.\nIt is an even further aim of the present invention to eliminate the transmission of disease from one patient to the next without having to dispose of the entire handpiece following every laser procedure.\nBriefly stated, the present invention describes a multi-component handpiece that is both economical and practical for surgical laser treatment. A disposable optical fiber is inserted into a reusable handpiece. The distal end of the optical fiber is protected by a micro-walled protective tube to prevent the disposable fiber from chipping before and during insertion into the handpiece. Once the optical fiber is protected, it is inserted into the handpiece by threading the optical fiber through the cap, body, and needle cannula of the handpiece until the fiber extends slightly beyond the distal end of the cannula. The micro-walled protective tube is removed from the end of the optical fiber, the optical fiber is positioned, and the cap is tightened. The cap and body of the handpiece cooperate to produce a tight friction fit that prevents longitudinal movement of the optical fiber. The tight friction fit holds the optical fiber in place without the use of adhesives to eliminate the problem of residual glue on the handpiece device and to allow for facilitated removal of the disposable fiber after use. Furthermore, a disposable fiber allows for convenient resterilization of the handpiece, eliminating the transmission of disease from one patient to the next without having to dispose of the entire handpiece following every laser procedure.\nThe above and other objects, features and advantages of the present invention will become apparent from the following detailed description read in conjunction with the accompanying drawings."}
{"text": "The present invention relates to monomers that are useful to make polybenzazole polymers, processes to make those monomers and methods to use those monomers.\nPolybenzazole polymers are synthesized by the reaction of AA-PBZ monomers with BB-PBZ monomers. The AA-PBZ monomer (which is sometimes also referred to as a dicarboxylic acid monomer) contains two electron-deficient carbon groups linked by a divalent organic moiety. It is frequently illustrated by Formula 1: EQU Q-DM-Q (1)\nwherein each Q is an electron-deficient carbon group and DM is a divalent organic moiety.\nThe BB-PBZ monomer (which is sometimes also referred to as a diamino-dihydroxy monomer or a diamino-dithio monomer) contains an aromatic group having two o-amino-basic moieties bonded to it. Each o-amino-basic moiety contains a first primary amine group bonded to the aromatic group and a hydroxyl, thiol or primary or secondary amine group bonded to the aromatic group in the ortho position with respect to the first primary amine group. The BB-PBZ monomer is ordinarily represented by Formula 2: ##STR1## wherein Ar.sup.1 is the aromatic group, each Z is --O--, --S--, or --NR--, and R is a hydrogen atom, an alkyl group or an aromatic group. The representation of bonds to an aromatic group without indicating position, as in Formula 2, indicates that the amine groups and --ZH groups may be in either cis- or trans- position with respect to each other. Each amine group must be ortho to one of the --ZH groups.\nThe BB-PBZ monomer is highly susceptible to oxidation when it is in a free-base state. It is ordinarily stored and handled as a hydrogen chloride salt, which is more stable. The hydrogen chloride is removed in a devolatilization step after the BB-PBZ monomer is dissolved in a polymerization solvent but prior to polymerization. The devolatilization step is time-consuming, but necessary because the hydrogen chloride interferes with polymerization of the monomer. It would be desirable to find new monomers that eliminate the need for the devolatilization step.\nThe monomers are usually polymerized in a dehydrating and non-oxidizing acid mixture. Proper stoichiometry is very important in the polymerization, since the reaction is a condensation polymerization. A small excess of BB-PBZ monomer can prevent the reaction from producing high-molecular weight polybenzazole polymer. It would be desirable to find new monomers that are simple to load in stoichiometric quantities."}
{"text": "Hydraulic motors or pumps currently available are of the type with radial or axial pistons and generally consist of a piston body, which acts on a rotatory eccentric shaft, which is inner in the case of the radial shape or rear in the case of the axial shape, whether the pistons of the piston body are straight or inclined. Said machines are adapted to convert the alternating motion of the pistons into the rotatory motion of the shaft with very high operating torques and pressure and can be used both as pumps and as motors in several activities of the construction industry. They all present the drawback that no equipment or apparatus of other associated machines is allowed to be arranged inside of them or pass through them. In these applications, generally on ordinary or planetary reduction gears, it is necessary to move the motor or pump to a position that is arranged on the side of the axis along which the torque is used, and to place it on a specially provided reduction gear, thus causing a consequent increase in the space needed and in manufacturing costs as well as a dispersion of power and, therefore, a decrease in mechanical efficiency. The object of the present invention is to provide a machine, in which the actuating pistons are arranged in a position that is tangential with respect to the rotation axis of the shaft that transmits the torque, thus allowing the creation of any kind of inner free space for the uses needed. Said pistons act on shafts that are arranged in a planet position with respect to the driving shaft and, in turn, transmit said torque by means of gearings, which can be of the ordinary or planetary type. The following description is relative to a motor with a number of 3 planets, but it also equally applies to a higher number of planet positions."}
{"text": "Most high bandwidth video acquisition systems record data for only short periods of time and are incapable of acquiring continuous high resolution imagery. Typically and more particularly, standard commercial video systems operate by approximately a factor of 10 less than the high resolution video acquisition system. Current recording systems have no capability to capture high resolution imagery and simultaneously record data for long periods of time (i.e., hours in lieu of minutes). Additionally, established high bandwidth digital scan converters produce analog outputs exclusively."}
{"text": "The HomePlug 1.0 standard provides for high-speed communications on a wired medium (e.g., powerlines and coaxial cables) using orthogonal frequency division multiplexed (OFDM) signals and carrier-sense multiple access with collision avoidance (CSMA/CA). HomePlug Audio/Video (HPAV) is a next generation of this powerline networking technology that provides for broadband high-quality multi-stream entertainment-oriented networking over a wired medium. HPAV networks use a central coordinator (CCo) to coordinate activities within an HPAV network including allocating time for CSMA and time-division multiple access (TDMA) transmissions.\nThe ITU-T G.9960 standard, commonly referred to as G.hn, also provides for broadband communications over a wired medium using OFDM signals. One issue is that HomePlug nodes and G.hn network nodes operating on the same wired medium may interfere with each other. To prevent G.hn network nodes and HomePlug nodes from interfering with each other, a frequency division multiplexing technique has been conventionally used to restrict signals from G.hn network nodes from the frequency band utilized by the HomePlug nodes. The HomePlug nodes however are configured to utilize a 2-28 MHZ frequency band, which is considered a lower-attenuation and lower noise region of the spectrum that is available on a powerline channel. As a result, G.hn network nodes are sometimes limited to using higher-attenuation and higher-noise regions of the available spectrum resulting in reduced performance for a G.hn network.\nThus, there are general needs for a methods and devices that provide a coexistence mechanism for HomePlug networks and G.hn networks to allow G.hn network nodes and HomePlug nodes to operate within the same frequency spectrum."}
{"text": "The present invention relates generally to the field of chemical treatment of water used for hydrostatic tests, and, more particularly, to a method and chemistry for the reduction of oxygen and microbiologically influenced corrosion in systems requiring a hydrostatic or leak test.\nWater used for hydrostatic or leak testing is typically left in pipelines from a couple of days to years. Thus, a pipeline undergoing a hydrostatic test may become quite vulnerable to oxygen corrosion and microbiologically influenced corrosion (MIC). MIC is often established during preparational phases of plant construction or during hydrostatic testing phases, when water is first introduced into the system. Therefore, this period is critical for prevention of oxygen corrosion and MIC and therefore, the quality of the water used for hydrostatic test is of vital importance.\nIn relation to a hydrostatic test, the main factors that cause corrosion of steel, leading to pipeline damage, are dissolved oxygen and bacteria present in the water used for the hydrostatic test. To prevent oxygen corrosion, oxygen scavengers such as ammonium or sodium bisulfites are commonly used. Preventive measures for control of MIC include biocide treatment, pH adjustment, or sulfate ion removal.\nThe use of biocides has proven to be effective in controlling MIC. However, most biocides add toxicity to the treated water. The toxicity of the water creates a problem because after the completion of a hydrostatic test the water must be discharged from the tested system, often with severe environmental constraints. For example, if the water is discharged into the ocean under U.S. jurisdiction, hydrostatic test water must pass certain aquatic toxicity tests. Due to such environmental constraints, the chemical selected to treat hydrotest water must not only inhibit oxygen corrosion and MIC, but must also comply with environmental requirements when discharged.\nThe adjustment of pH inhibits the bacterial growth to a certain extent. Further, the mere adjustment of the pH of water does not protect pipe surfaces from oxygen corrosion, while potentially causing mineral scale problems. Thus, there remains a need for a chemical treatment regimen for water used in hydrostatic test procedures that reduces or prohibits oxygen and microbiologically influenced corrosion of ferrous metals while also meeting environmental constraints placed on the discharge of such water. The present invention is directed to solving this problem in the art.\nThe present invention combines a chemical treatment strategy, including biocide, oxygen scavengers, and a scale inhibitor, with an adjustment of pH. This combination chemical stratagem provides an effective option for treating hydrotest water for inhibition of oxygen corrosion, inhibition of MIC, and the safe discharge of water into the environment following the performance of a hydrostatic test.\nThe primary benefit of the combination treatment plan is reduced biocide usage in a high ( greater than 9.0) pH brine. The reduced overall chemical usage facilitates meeting environmental guidelines while also minimizing chemical cost. The method of the invention therefore comprises adding an oxygen scavenger to remove oxygen and prevent or minimize oxygen corrosion; raising the hydrotest water pH, typically in excess of pH 9.5 (this may be accomplished with many bases, although sodium hydroxide is typical); adding a biocide in reduced amount (the ability of the biocide to function in reduced quantity is allowed by the increase pH of hydrotest brine); and adding a scale inhibitor to inhibit scale (an increase in pH of the hydrotest brine can increase the tendency to form mineral scales, such as calcium or magnesium carbonate).\nThe order of the addition of the chemicals to the hydrotest water may be important, particularly if using oxygen scavengers such as bisulfite that can react with and impair the performance of many biocides. If the selected oxygen scanvenger adversely reacts with the biocide, the it is important to add oxygen scavenger prior to adding biocide and to allow adequate time for reaction.\nBased on the previous scenario, a typical application would be, to first inject a package consisting of sodium hydroxide solution, oxygen scavenger, and scale inhibitor into hydrotest water. This would then be followed by injection of biocide and used for the hydrotest application. If there are no negative interactions between the selected chemicals, then all of the chemicals may be injected simultaneously or in a single package.\nThese and other benefits, features, and advantages of the invention will be apparent to those skilled in the art from a review of the following detailed description along with the accompanying drawings."}
{"text": "The present disclosure relates to hydrostatic regenerative braking for efficient use of energy of a vehicle. More particularly, a transmission for coordinating the flow of power between a prime mover, a hydraulic energy storage system and the vehicle wheels is discussed.\nRegenerative drive systems including hybrid hydraulic arrangements have been applied to motor vehicles in the past. While conventional vehicle braking systems typically convert a vehicle's kinetic energy into heat energy, hydrostatic regenerative braking systems convert a moving vehicle's kinetic energy into stored hydraulic energy. The hydraulic energy is typically stored in an accumulator for later use to propel the vehicle. At least one hybrid hydraulic drive system includes a hydraulic pump/motor selectively operable to transfer energy from the vehicle driveline to a hydraulic storage device, such as during a braking event, or transfer energy from the hydraulic storage device to the driveline, such as during vehicle acceleration. When the hybrid hydraulic system stores power from the driveline, the hydraulic pump/motor acts as a pump to provide pressurized fluid to the accumulator. When the hybrid hydraulic system transfers power to the driveline, the hydraulic pump/motor acts as a motor driven by the energy stored in the accumulator.\nSome vehicles include transmissions for transferring power between the driveline and the hydraulic storage device. At least one transmission cooperates with a hydrostatic regenerative braking system and includes a two gear design with a clutch positioned at an output shaft of the transmission. While this design may provide some benefit, it requires a relatively large diameter driven gear to provide a useful gear reduction ratio. Large packaging volume requirements and high weight are associated with this transmission. Due to the clutch being positioned downstream of the gear reduction mechanism, the clutch is required to transmit a torque approximately three times greater than the pump torque. A relatively large clutch is required to transfer the torque. The size and weight of the transmission are accordingly increased. Furthermore, with the two gear design, only one output drive ratio is available. To provide different versions with different ratios, new sets of gears for each ratio must be provided. Accordingly, a need exists in the art for an improved hydrostatic regenerative braking transmission."}
{"text": "1. Field of the Invention\nThe present invention generally relates to microprocessor and multiprocessor architectures and, more particularly, to thread-level speculative processor execution for achieving high performance and speeding up sequential applications.\n2. Description of the Prior Art\nAs increasing numbers of smaller and faster transistors can be integrated on a single chip, new processors are designed to use these transistors effectively to increase performance. The arising challenge is to find the most effective way to put these transistors in use. Currently, most computer designers opt to use increasing transistor budget to build even bigger and more complex uniprocessors. Another possibility is to place big amounts of memory on the chip. Alternatively, multiple processor cores can be placed on a single chip. The later approach is called chip multiprocessors (CMP).\nPerformance improvements using a single complex processor is achieved by exploiting ILP (instruction level parallelism), i.e. by finding non-dependent instructions in a program sequence which are then executed at the same time. However, the possible performance gain by exploiting IPL is limited due to finite amount of ILP present in any particular application sequence.\nPlacing multiple smaller processor cores on a single chip is attractive because single processor core is less complex to design and verify. This results in less costly and complex verification process as once verified module-processor—is repeated multiple times on a chip. Each processor core on a multiprocessor is much smaller than a competitive uniprocessor, minimizing the core design time. In addition, keeping design partitions small—like a single processor core in a CMP—design tools can handle processor complexity much easier, compared to competitive complex uniprocessors. However, many important existing applications are written for uniprocessors, and it is a non-trivial task to convert uniprocessor applications into multiprocessor ones. For this, sequential programs have to be explicitly broken into threads and synchronized properly. So far, parallelizing compilers have been only partly successful at automatically handling these tasks.\nSpeculative multithreaded processors present possible solution of these difficulties offering high potential performance improvement. A speculative multithreaded processor consists logically of replicated processor cores that cooperatively perform the parallel execution of a sequential program. The sequential program is divided into chunks called speculative threads, and these threads are executed on processor cores concurrently and speculatively. This approach for performance improvement by exploiting coarse-grain parallelism in addition or instead of fine-grain parallelism (e.g., ILP) is called thread level speculation (TLS). In thread level speculation approach, sequential programs are divided into speculative threads which are then executed concurrently on processor cores. Ideally, there are no data and/or control dependences between the threads, but being parts of the same sequential program, speculative threads are both data and control dependant. The data flow between speculative threads in one direction only—from sequentially older threads to younger ones. (Thus, data used in a younger speculative thread can be a result calculated in an older thread.) To ensure that each program executes the same way that it did on a uniprocessor, hardware must track all inherited dependences. When a younger thread in a sequence causes a true dependence violation, the hardware must ensure that the misspeculation is detected, and the misspeculated thread has to re-execute with the correct data.\nTo support speculation, multiprocessor architecture for thread level speculation has to fulfill the following requirements: 1) it has to maintain a notion of the relative order of the threads—i.e., know which thread is executed before some other thread in a sequential program; 2) it has to forward data between parallel threads, or predict data; 3) it has to support mechanism for dependency violation detection—to detect if read occurred too early; 4) it has to safely discard speculative thread once dependency violation is detected; 5) it has to commit speculative writes in proper order—only after making sure that this thread would have been executed the same in a sequential execution; and, 6) it has to re-execute the misspeculated threads with proper data.\nA number of multiprocessor architectures with support for thread level speculation have been proposed. In several of these architectures, a program is chopped into threads by the compiler during the compilation time, such as in a multiscalar processor as proposed in the reference to G. S. Sohi, et al. entitled “Multiscalar Processors”, 27th International Symposium on Computer Architecture (ISCA-22), 1995, or as in a superthreaded architecture or trace processor. In other approaches, hardware dynamically forms the threads during the run time, such as proposed in the reference entitled “Dynamic Multithreaded Processor” by H. Akkary and M. Driscoll in Proc. Of the 31st Annual International Symposium on Microarchitecture (1998) and “Clustered Speculative Multithreaded Processor” proposed by P. Marcuello and A. Gonzales in Proc. Of the 13th Intl. Conference on Supercomputing, pp. 365-372 (1999). All of these architectures require significant changes on the processor core or/and on the L1 and/or L2 level caches to support thread level speculation. These changes include at least one of the following: 1) provision of means for registers forwarding between processors; 2) the addition of new fields in one or more caches to distinguish speculative vs. non-speculative values; 3) a modified processor interface to allow communication of speculative values; and 4) a change of speculation status for the processor. Requiring significant changes to the processor core and/or to the memory nest to enable thread level speculation, existing architectures can not take advantage of increased performance which TLS offers. To support thread level speculation on the existing processor, the processor core needs massive re-design and complete re-verification process. Similarly for the memory nest, re-design and verification effort makes it prohibitive, or at least very expensive, for already existing cores and system.\nIt would be highly desirable to provide a system and method which would enable thread level speculative execution on existing processors and memory systems without requiring costly changes to the processor core or existing cache hierarchy."}
{"text": "In the field of communications, the need for high-speed transmission of data including video and audio has continued to increase. Moreover, there has been an increase in the selection of services by which users can connect to a network, such as the Internet. Specifically, Internet Service Providers (ISPs) may allow for connectivity to the Internet through lower-speed connections at different rates, such as 56 kilobits/second, by employing a Plain Old Telephone Service (POTS) line. Other choices for connection, which are at higher speeds, into a network can include Integrated Services Digital Network (ISDN), Digital Subscriber Line (DSL) service (both using a POTS line), and cable modem service over a Radio Frequency (RF) cable line.\nAdditionally, there is an increasing demand by the ISPs to allow for modifications in the connection rates for a given type of service. In particular, the market to provide certain services, including DSL, is very competitive. Accordingly, the low-speed connections are deeply discounted. Additionally, there are less users of the higher-speed connections due to the costs. Therefore, ISPs need an approach to offering variable services, wherein a subscriber can typically pay for low-speed connection service, while occasionally upgrading their connection when a higher bandwidth may be needed and/or advantageous.\nAccordingly, this enables the ISPs as well as the users of the ISPs flexibility with regard to bandwidth usage and costs associated therewith. To help illustrate, a given user may want to connect to the Internet at a slower connection speed when the user is only traversing pages of different web sites, in comparison to when the user may be involved in a video conference and/or downloading a large amount of data wherein a higher connection speed is needed and/or more advantageous.\nDisadvantageously, current approaches to allow for modifications to a subscriber connection (such as modifications to the connection rate or policing values for the connection) to a network require that the subscriber connection be broken, wherein a new connection is re-established that applies these modifications. Accordingly, subscribers experience service interruptions when a modification is made to the current connection. Moreover, administrators working for the ISPs can be required to manage the updates to these subscriber accounts and/or connections."}
{"text": "In memory circuits utilizing variable threshold devices, such as MNOS transitors, the threshold value at which the transistor conducts may be controlled to have a low threshold level for supplying a logic or binary \"one\"; and to have a high threshold level fro supplying a logic or binary \"zero\", for example. In such circuits, it is desirable to have a short write cycle time to enable fast writing of data into the memory cell as well as reading such data out. In a prior art MNOS memory circuit, data may be written into the memory cell over and over. In other words, the memory cell may be written as a binary one for a number of times in a row before a binary zero is written. This causes the variable threshold transistor to shift its threshold between high and low to the maximum extent during the write cycle which saturates the device to a condition. which may be referred to as the saturated threshold state. Thus, when it is desired to write a binary level of the opposite state into the memory cell, the write pulse must have adequate polarizing voltage and time to shift the threshold voltage of the transistor from the previous saturated state. To speed up the time of writing a binary level or shifting the voltage in a variable threshold transistor, such as an MNOS transistor, the voltage or electric field across the gate insulator of the device is increased by increasing the polarization voltage. Unfortunately, these high electric fields across the gate insulator of the MNOS device during the write cycle serve to accelerate the undesirable endurance phenomena creating a smaller threshold voltage and decreased retention time of the memory state.\nIn an attempt to prevent saturation of the transistor, one typical prior art circuit operated such that the transistor was cleared prior to each writing of the transistor to a high threshold state. Such clearing consisted of applying a voltage to obtain a low threshold state prior to each writing of the data, regardless of the data already in the memory cell. This was followed by applying the voltage for the high threshold state to selected memory cells. Such a system, which required the application of a write voltage first for logical zeroes, and then for logical ones, required more than twice the minimum write time.\nIn order to overcome the disadvantages of operation created by the saturated threshold state without increasing the polarization voltage or without requiring that each cell be operated to its low threshold state prior to selectively operating the memory cells to the high threshold state during each write cycle, a system and method described in U.S. Pat. No. 4,090,258, dated May 16, 1978 and issued to J. R. Cricchi, one of the inventors of the present application, is proposed. This system describes a variable threshold MNOS non-volatile memory system and method that utilizes a pair of MNOS transistors for each memory cell. For a binary one, one transistor is operated to its high threshold state, while the other transistor is simultaneously operated to its low threshold state. For a binary zero, the one transistor is operated to its low threshold state while the other transistor is operated to its high threshold state. Prior to each write cycle of the circuit, a read-compare operation is performed, which compares the data to be written with the data actually stored in the memory cell. In the event that the data stored in the memory cell is the same as the data to be written into that memory cell, the write portion of the cycle is inhibited. This inhibiting of the rewriting of the same data into the memory cell prevents the transistor from being operated to its saturation threshold state without the necessity of clearing the transistor prior to the writing in of each new data. Further, in such system the fact that one of the transistors of the cell is operated to its low threshold state and the other to its high threshold state increases the \"window\" to minimize the errors when reading out the data stored in the memory cell.\nHowever, MNOS memories that have fast write characteristics tend to have short retention periods. Although the proposed system provides for a large window, it is possible that certain memory cell remain in their binary one or binary zero state during extended operations that the detection window would decrease; that is, the high threshold state of one of the transistors of the memory cell would decrease and the low threshold state of the other transistor of the memory cell would increase. Although such a system provided accuracy of operation down to a one-half volt difference between the two transistors of a memory cell, it is desirable that the transistors be refreshed periodically without driving the transistors to their saturated threshold state. Also, it is desirable that the refreshing of the memory cells be controlled internal to the memory array rather than by external programming."}
{"text": "Internal combustion engines such as diesel engines, gasoline engines, and gaseous fuel-powered engines exhaust a complex mixture of air pollutants as byproducts of the combustion process. These air pollutants are composed of gaseous compounds including, among other things, the oxides of nitrogen (NOx). Due to increased attention on the environment, exhaust emission standards have become more stringent and the amount of NOx emitted to the atmosphere from an engine can be regulated depending on the type of engine, size of engine, and/or class of engine.\nOne method that has been utilized by engine manufacturers to comply with the regulation of engine emissions has been to implement exhaust gas regeneration (EGR). EGR systems recirculate exhaust gas to mix with combustion air supplied to an engine. The recirculated exhaust gas reduces a concentration of oxygen and increases a thermal mass within the engine's cylinders, thereby lowering a resulting combustion temperature. The lowered combustion temperature slows the chemical reaction of the combustion process and decreases the formation of NOx.\nAlthough successful at reducing NOx, the effectiveness of an EGR system can be influenced by an amount of heat contained within the exhaust gas that is recirculated. That is, the exhaust gas that is recirculated back into the engine can contain large amounts of energy in the form of waste heat, which can increase the in-cylinder combustion temperature and counteract the desired NOx reduction. In addition, the excess heat also taxes the engine's cooling system and may represent an unutilized source of energy that, if harnessed, could increase engine efficiency.\nAn exemplary system implementing EGR is disclosed in U.S. Pat. No. 6,216,458 of Alger et al. that issued on Apr. 17, 2001 (“the '458 patent”). Specifically, the '458 patent discloses an exhaust gas recirculating (EGR) system for a turbocharged internal combustion engine. The EGR system includes a recirculation conduit for recirculating a volume of exhaust gas from the exhaust manifold to the intake manifold, a cooling air conduit for transporting a flow of cooling air, and a gas-to-air heat exchanger disposed in operative association with the recirculation conduit and the cooling air conduit. The heat exchanger is adapted for cooling the volume of the recirculated exhaust gas in the recirculation conduit."}
{"text": "Probe data may be collected by various data collection devices including global positioning systems (GPS) systems and other navigation systems. Regardless of the device that collects the probe data, the probe data may include a plurality of probe points. Each probe point may identify the location, such as in terms of latitude and longitude, at which the device that collected the probe data was located at the time at which the probe point was captured. A probe point may also include or otherwise be associated with a z-axis component indicative of the altitude or height of the device at the time at which the probe point was captured. However, information relating to the altitude associated with the probe point generally has a much greater tolerance and a much larger error associated therewith than the latitude and longitude associated with the probe point. For example, the altitude associated with a probe point may have an error of 10, 20, 30, 40, 50 or more meters. Additionally, not all probe points have an associated altitude with some probe points defining a location merely in terms of latitude and longitude.\nDuring the creation of maps that identify road segments, the error associated with the altitude of the probe point has made it particularly difficult to determine if the probe points are associated with a road that passes over another road or a road that passes under another road. In this regard, roads sometimes cross one another or are otherwise stacked in such a manner that one road passes beneath another road. For example, an interchange of an interstate may include an underpass and an overpass with the road that forms the underpass passing beneath the road that forms the overpass. By way of another example, roads may be stacked with one road on top of another even in instances in which the roads do not cross. In these instances, the error associated with the altitude of the probe points may be sufficiently large that the probe points may not be reliably associated with either the upper or the lower road with a sufficient degree of confidence in the map creation or refinement process. In this regard, the error associated with the altitude of the probe points may be so large as to equal or exceed the anticipated vertical spacing between the roads such that the probe points may not be matched to the road that forms the overpass or the road that forms the underpass with the desired degree of confidence."}
{"text": "The present invention relates to a heavy-duty vehicle combination or a heavy-duty vehicles compound, such as for example tractors train, trailer power vehicles, trailer train or platforms train, particularly with turnable flat bed or platform vehicles supported by a plurality of axles or wheel groups and with at least partially driven axles or wheel groups by an additional drive on the load carrying vehicles.\nHeavy-duty vehicle systems of the above mentioned general type are known in the art. With heavy loads, particularly loads which must be transported with consideration of special regulations or particular permission on street vehicles, it is in general not effective to transport the working load exclusively on pulled or pushed vehicles, trailers, etc., inasmuch as a tractor or a pushing vehicle which is the only driven vehicle in system must be loaded with high ballast weight in order to guarantee traction and friction loads required for movement of the vehicle combination.\nThe German patent 841,850 discloses trailer vehicles whose wheels are driven hydraulically by a motor unit arranged on a tractor, which in accordance with the power requirements can be turned on individually or in a parallel operation. In this construction it is not possible to change the ratio of the drive power supplied to the individual vehicles of the vehicle compound."}
{"text": "(1) Field of the Invention\nThe present invention concerns a system for controlling hydraulic pressure in an automatic four-speed transmission, and more specifically a means for stabilizing line pressure and control pressure as well as reducing the shifting shock when making a gear shift from the neutral mode to the forward or reverse drive mode.\n(2) Description of the Prior Art\nThe automatic transmission comprises a torque converter, multi-stage gear shift mechanism connected with the torque converter, and a plurality of friction elements for selecting one of the gear shift stages. The friction elements are worked through a plurality of control valves for controlling the hydraulic pressure applied by the hydraulic control system.\nSuch a hydraulic control system comprises a pressure regulating means for regulating the hydraulic pressure produced by an oil pump, manual and automatic transmission control means for selecting one of the shift modes, pressure control means for controlling the shifting sense and responsiveness to make a gear shift smooth, damper clutch control means for working the damper clutch of the torque converter, and pressure distribution means for distributing properly the hydraulic pressure among the friction elements. The pressure control means controls the pressures applied to the friction elements, torque converter, solenoid valves, etc., which substantially affects the shifting sense.\nWhen performing the shift operation, it is preferable to supply controlled hydraulic pressures to the friction elements while completing the shift operation with the drive pressure. Such transmission control means is disclosed in U.S. patent application Ser. No. 08/347,312 filed by the present applicant, which suffers the shift shock caused by the drive pressure or the backward pressure applied to the friction elements in the manual shifting when making the forward or the backward movement."}
{"text": "1. Field of the Invention\nThis invention relates to computerized record processing systems and more particularly to systems and methods for efficiently processing a collection of records through one or more decision trees.\n2. Background of the Invention\nThe computation time required for certain types of record processing increases rapidly as the number of records increases. For example, record linkage requires comparing pairs of records. Each such comparison is computationally expensive. Additionally, as the number records increases, the number of comparisons that need to be conducted grows exponentially. Accordingly, what is needed is a computer system configured to efficiently process large numbers of records."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an Orthogonal Frequency Division Multiple Access (OFDMA) system. In particular, to the present invention relates a method and apparatus for detecting a cell in an OFDMA system.\n2. Description of the Related Art\nMobile communication systems have been evolving into a 4th generation (4G) mobile communication system providing super high-speed multimedia service, succeeding a 1st generation (1G) analog system, a 2nd generation (2G) digital system and a 3rd generation (3G) IMT-2000 system providing high-speed multimedia service. In the 4G mobile communication system, a subscriber can access all of a satellite network, a local area network (LAN), and an Internet Protocol (IP) network using one mobile station (MS). That is, the subscriber can enjoy voice, image, multimedia, Internet data, voice mail, and instant message (IM) services using one mobile station. The 4G mobile communication system supports a data rate of 20 Mbps for super high-speed multimedia service, and uses an orthogonal frequency like an Orthogonal Frequency Division Multiplexing (OFDM) scheme.\nThe OFDM scheme, a digital modulation scheme for multiplexing a plurality of orthogonal carrier signals, divides a single data stream into several low-rate streams and simultaneously transmits the low-rate streams using several low-rate subcarriers. As a result, a symbol interval increases, causing a reduction in relative dispersion in a time domain due to multipath delay spread.\nAn Orthogonal Frequency Division Multiple Access (OFDMA) system transmits data per symbol. Interference occurs between the symbols, and in order to compensate for the intersymbol interference, the OFDMA system inserts a cyclic prefix (CP) which is longer than a transport channel, into the symbol.\nFIG. 1 is a diagram illustrating a symbol structure in an OFDMA system. Referring to FIG. 1, oblique-lined regions correspond to the CP. A rear part of the symbol is copied and then attached to a front part for a given guide time Tg. Here, a time defined by excluding the CP from the symbol is denoted by Tb, and a time comprising the full symbol is denoted by Ts.\nIf the number of subcarriers used is denoted by N, a reception signal obtained after CP removing and fast Fourier transform (FFT) has the following relation.z(k)=H(k)s(k)+ω(k)   (1)\nIn Equation (1), s(k) denotes a reception signal in a frequency domain, H(k) denotes a value obtained by performing N-point Discrete Fourier Transform (DFT) on a time-domain channel response h[n], and ω(k) denotes an N-point DFT coefficient for a white Gaussian noise ω[n] and has a dispersion of N0. Herein, [n] and (k) are factors for representing a time-domain signal and a frequency-domain signal, respectively.\nA mobile station is required to estimate a channel H(k) in order to demodulate a signal received from a base station (BS), and to this end, the base station inserts pilots in a downlink data packet before transmission. With the use of the pilots, the mobile station not only performs channel estimation but also estimates signal-to-interference and noise ratio (SINR) information useful for power control in a multiple access scheme and transmits the SINR information to the base station.\nIn a cellular system, a mobile station should detect a cell to which it belongs in order to initiate communication. The cell detection is achieved through unique pseudo-random noise (PN) codes used in each base station and a cross-correlation between mobile station's reception signals. In a Wideband Code Division Multiple Access (WCDMA) system, the cell detection is achieved through a PN code allocated to each cell and a cross-correlation between a primary synchronization channel (P-SCH), a secondary synchronization channel (S-SCH) and a common pilot channel (CPICH) at the beginning of communication. On the other hand, in the OFDMA system, a base station sends a PN code allocated thereto using a preamble inserted at the head of a data frame and a mobile station can detect a cell through cross-correlation. However, the cross-correlation needs N2 multiplications, causing an increase in the number of calculations.\nFor the frequency-domain data signal of Equation (1), the time-domain reception signal z(k) that a mobile station receives can be represented by the product of a channel frequency response and a frequency-domain transmission signal as expressed in Equation (2).z[n]=h[n]⊙Ns[n]+w[n]  (2)where ⊙N· denotes N-circular convolution, h[n] denotes a channel response in the time domain, and ω[n] denotes a white Gaussian noise in the time domain.\nBased on the relation of Equation (1), equalization or channel estimation is achieved through N efficient divisions. For equalization, particular data s(k) is estimated by dividing z(k) by an estimated value {tilde over (H)}(k) for H(k). However, for cell detection, s(k) is limited to several particular PN codes and a PN code allocated to a corresponding cell among the PN codes should be detected without an estimated value for H(k). In this case where there is no condition for H(k), it is impossible to detect a cell. However, in the OFDMA system, it is possible to detect a cell on condition that a channel length L in the time domain is very smaller than the number N of OFDM symbol sampling. In general system implementation, however, although information on the channel length L should be given to a mobile station, it is difficult for the mobile station to acquire such information at the beginning of communication."}
{"text": "One of the most common light fixtures is the recessed can downlight (RCD), which is an open-bottom can that contains a lightbulb, most commonly an incandescent bulb or a fluorescent bulb. The fixture is typically connected to the power mains at 120 to 277 volts, 50/60 Hz. RCDs are generally installed during the construction of a building before the ceiling material (such as plaster or gypsum board) is applied. Therefore, they are not easily removed or substantially reconfigured during their lifetime.\nRCDs generally also accommodate lightbulbs having various sizes, different overall dimensions (i.e., length, width, and diameter), and varied light-distribution capabilities. For example, various bulbs have narrow, medium, or wide (flood) light distributions. Therefore, the internal features of the RCD are constructed to accommodate many (if not all) different bulb types. Such features include mechanisms to adjust the vertical position of the bulb socket, as well as reflectors that channel and distribute the light. Because there are so many different lightbulbs and finishes, a very large number of trim rings and optics combinations may be utilized in RCDs, in addition to the various spacers that accommodate the bulbs. Thus a complex arrangement of parts is needed for each RCD that is produced.\nBecause LEDs have very high efficiency (e.g., 100 lumens per watt compared to 10-15 lumens per watt for incandescent or halogen lights) and a long lifetime (e.g., 10,000-100,000 hours), they are attractive for virtually all lighting applications. LED retrofit fixtures have been designed to replace existing, installed RCD fixtures. U.S. Ser. No. 14/660,159, filed on Mar. 17, 2015, for example, describes a retrofit kit that enables retrofitting of a wide variety of different RCDs (e.g., RCDs incorporating fluorescent bulbs) with a single “universal” LED-based fixture that is quickly and efficiently installable. Within the retrofit kit, the LED light sources and control electronics are modularized for ease of assembly and installation. In addition, the retrofit kit may be utilized substantially independently of the specific lightbulb being replaced yet conforms to the volume and desired level of illumination of the existing RCD.\nA retrofit kit as described in the '801 application may include a discrete driver module featuring circuitry for supplying power to and controlling the LED light source(s), as well as, in various embodiments, circuitry for controlling the LEDs based on sensed temperature (for example, the temperature of the LEDs themselves or of one or more temperature sensors such as thermistors in close proximity to the LEDs). The driver module is electrically connected to a discrete lighting module featuring one or more LEDs (for example, several LEDs arranged in a rectilinear array) via a flexible conduit that contains and protects one or more wires carrying electrical signals between the two modules. The lighting module may incorporate one or more temperature sensors for sensing the temperature of the LED(s) and/or the ambient temperature, and the driver module may incorporate thermal-feedback circuitry for controlling power supply to the LED(s) based on the sensed temperature. The lighting module also typically incorporates an integral or removable heat sink.\nUnfortunately, the high power levels often required to drive an LED retrofit solution to maintain previous lighting levels may generate so much heat that merely heat-sinking the LEDs can prove insufficient. LED lifetime can be substantially shortened by excessive operating temperatures; in general, it is advisable to maintain the LED below 100° C. during operation. Indeed, even where such passive measures as finned heat sinks are sufficient from a performance perspective, they may be incompatible with the physical restrictions of a retrofit; the volume within a light source such as an RCD is limited, and the airflow needed for effective heat sinking may be impossible within the fixture space. Even when there is adequate room for a large heat sink, it may displace the light source so as to create glare and ultimately impose a cap on light output. For example, it may be necessary to change the configuration to position the LEDs lower in the can, resulting in an out-of-focus condition for the LEDs and/or considerable visual glare, which is highly undesirable. Placing the heat sink outside the can (reflector) also is usually not possible due to the mounting and support structure of the light fixture. Moreover, the region above the can may be filled with insulation and the building's structural elements, such as rafters and beams, may either restrict airflow or make the use of the space impossible."}
{"text": "Herbal extracts and formulations have been suggested for the trichological treatment of humans, particularly males, but not exclusively so. Efficacy varies from nil to some reduction in hair loss and/or noticeable hair regrowth on the heads of patients. However, the interaction of the constituent compounds in such formulations with one another is little understood. The respective modes of action of such formulations are quite diverse. Generally, such treatments are either topical or systemic.\nThe term alopecia is used in the art generally to refer to all conditions of baldness, including partial baldness, and the term alopecia is used in this specification to have the same meaning, as well as to refer to hair loss generally, unless the context requires otherwise. The term “trichological treatment” is used to refer to the treatment and/or formulations for the treatment of alopecia and hair thinning in males and females. Although the term is used sometimes to refer to scalp conditions or the like, the invention is not limited to treatment of such conditions, and may be used to treat hair loss on other parts of the body.\nThere are several root causes of alopecia, some of which are well established both clinically and pathologically, and quite well understood. Other root causes are not well understood and/or subject of conjecture. Some herbal and drug treatments of alopecia are directed to one or more of the root causes of alopecia, whether they be well understood or merely conjectured. There is a quantity of research data on efficacy, some of which is in conflict, indicating that the alopecia condition is complex, and is probably affected by a range of different mechanisms. Some root causes are expressed or exhibited more than others, but unless several of such root causes are addressed by a proposed remedy, it is suggested that its efficacy may be limited.\nThe present invention aims to provide formulations and treatments for trichology which alleviate one or more of the shortcomings of the prior art. Other aims and advantages of the invention may become apparent from the following description."}
{"text": "A fuel vapor storage canister in an evaporative emission control system for a motor vehicle typically includes a mass of carbon granules in a carbon bed chamber of the canister which adsorb fuel vapor conducted to the canister from a fuel tank of the motor vehicle through a vapor transfer conduit. Fuel vapor is extracted from the carbon bed chamber through a vapor purge duct. Commonly, a liquid trap on the vapor storage canister captures liquid fuel entrained with the fuel vapor in the vapor transfer conduit to prevent contamination of the carbon granules. Liquid traps on some vapor storage canisters are passively purged, i.e., captured liquid fuel either evaporates from the trap or drains by gravity back to the fuel tank. Liquid traps on vapor storage canisters described in U.S. Pat. Nos. 4,714,485 and 4,853,009, issued Dec. 22, 1987 and Aug. 1, 1989, respectively, and assigned to the assignee of this invention, are actively purged, i.e., captured liquid fuel is induced by a pressure gradient between the liquid trap and the purge duct to flow from the liquid trap into the purge duct through a metering orifice submerged in liquid fuel in the liquid trap. A fuel vapor storage canister according to this invention is a novel alternative to the fuel vapor storage canisters described in the aforesaid U.S. Pat. Nos. 4,714,485 and 4,853,009."}
{"text": "Ultrasound technology in medicine continues to improve the ability of health care workers to obtain images of internal human structures. Not only are images obtained quickly, they also are delivered with minimal harm and discomfort to the patient. One current example of this technology is the real-time two-dimensional picture of the human fetus frequently obtained by expectant mothers when they get a sonogram by an obstetrician. These images are pie-shaped forms with the apex at the skin surface and showing ever deeper layers radiating out from the apex, namely, the wall of the uterus, the fetus, and then the opposite wall of the uterus. The depth of the image is limited by the frequency at which the sound is emitted, its amplitude and the currently available technology.\nWhile visualizing the fetus may be the best known use of this technology in medicine, it is also used to look at the skin, the liver, the ovaries, the kidneys, and the heart. In looking at the heart it is invaluable because it provides a functional study. By this is meant that a real-time moving image of the heart (as it contracts in front of the observer) is obtained, yielding information about the contractions. This knowledge is readily obtained from watching the moving image on a computer monitor. Some examples of the interesting information that can be obtained are how much each chamber of the heart pumps with each beat, and what percentage of blood is expelled from each chamber with each contraction.\nA simple explanation of ultrasound technology is as follows. An ultrasound transducer uses a piezoelectric crystal or similar device to convert electrical energy to sound vibrations and vice versa. High frequency sound energy is emitted by the transducer into a sound conductive medium. Since different structures in the human body conduct and reflect sound energy differently, the various times and intensity at which the sound wave energy is reflected and returned to the transducer can be used to reconstruct an image, using available computerized techniques. A simple analogy is the use of sonar to determine the range, bearing and size of a submarine located in the ocean. Sonar uses essentially the same principles of measuring reflected sound. In the study of the human body, the ranges to structures are much smaller, allowing hundreds of images to be obtained each second. This allows for the real-time visualization of moving structures.\nThe most commonly manifested use of ultrasound technology in medicine employs a phased array or other ultrasound transducer to obtain two-dimensional tomographic images of human structures. In other words, the images represent planar \"slices\" through the body. Different two-dimensional slices are obtained by adjusting the orientation and position of the transducer. An analogy would be the result obtained by a butcher when using a slicing machine to obtain thin sheets of ham. If the ham were rotated in the machine before each slice, different shaped slices could be obtained.\nFollowing in this thought is the fact that if the sequence of the slices of ham are known, a reasonable approximation of the original ham can be obtained by restacking the original slices. A similar premise lies behind presently available three-dimensional ultrasound imaging, wherein a series of two-dimensional images of a structure in the human body is used to recreate pictorially a three-dimensional image of that structure. This is not a new idea, with several approaches currently being summarized in patents and the scientific literature. The basic premise of the current art is that if the time of imaging and the location of the ultrasound transducer are known, the order and position of individual two-dimensional slices is also known. A relatively simple computer algorithm then \"stacks,\" or reassembles, these slices into a three-dimensional picture. The present application of this technology is limited by constraints in the accuracy of assumptions made about the position of the transducer as the individual two-dimensional images are taken. For example, it may be assumed that the position of the transducer is moving linearly, while, in fact, extraneous movements of the patient and the like may cause it to move non-linearly. Consequently, when the images are assembled (\"stacked\") according to such inaccurate positional assumptions, loss of focus and other image distortion results. In other cases, a constant velocity of movement is assumed, but this also is difficult to achieve in practice.\nSound waves are poorly transmitted by bone and air. This prevents adequate images from being obtained of, for example, the heart when attempting to look through ribs, the breastbone (sternum), or the lungs. In looking at the heart, there are two primary ways in which ultrasound beams can be directed. One way to avoid bone is by looking from the upper abdomen, just below the breastbone, and up at the heart. This is known as the transabdominal or subxiphoid view. The second way is by inserting a transducer probe into the esophagus--a structure which is in close proximity to the heart during much of its course from pharynx to stomach. This transesophageal view is widely regarded as allowing for superior ultrasound views of the heart. Different areas of the heart can be seen simply by sliding the transducer up and down in the esophagus. Obviously, this latter approach is less well tolerated by the patient, requiring sedation while the flexible probe (often 1 to 1.5 cm. in diameter) is manipulated in his or her esophagus. In addition, there are greater risks to the patient, including but not limited to, tears of the esophagus. The larger the probe and the longer it must remain in place, the greater are the discomfort and risks.\nPrior developments in this field may be generally illustrated by reference to the following information disclosure statement:\n______________________________________ U.S. PATENT DOCUMENTS U.S. Pat. No. Patentee Issue Date ______________________________________ 5,546,949 L. Frazin et al. Aug. 20, 1996 4,932,414 D. Coleman et al. Jun. 12, 1990 5,398,691 R. Martin et al. Mar. 21, 1995 5,295,486 H. Wollschlager et al. Mar. 22, 1994 4,798,210 R. Ledley Jan. 17, 1989 5,515,856 B. Olstad et al. May 14, 1996 5,396,890 L. Weng Mar. 14, 1995 ______________________________________"}
{"text": "1. Field of the Invention\nThis invention relates to a process for the hydrolysis of adiponitrile and production of nylon 6,6. More specifically, the invention relates to the catalytic hydrolysis of adiponitrile at low catalyst levels.\n2. Description of Related Art\nPolyamides are currently produced commercially from dicarboxylic acids and diamines. For example, nylon 6,6 is commercially produced from adipic acid and hexamethylenediamine. Alternate routes to these materials have been explored within the art. For example, the art has taught processes to produce nylon 6,6 from adiponitrile and hexamethylenediamine. U.S. Pat. No. 2,245,129 discloses a process to produce polyamides by reacting adiponitrile, hexamethylenediamine, and water at high temperatures and pressures. A modified process to react adiponitrile with hexamethylenediamine and water was disclosed in U.S. Pat. No. 3,847,876. Catalysts were introduced into the processes to promote this reaction as disclosed by U.S. Pat. Nos. 4,490,521; 4,542,205; 4,603,192; 4,725,666; 4,749,776; 4,436,898; and 4,528,362. U.S. Pat. No. 4,501,881 discloses a process to form polyamides from adiponitrile, hexamethylenediamine, adipic acid and water. A significant shortcoming of these before mentioned processes to produce polyamides from adiponitrile and hexamethylenediamine is the production of significant levels of the dimer of the aliphatic diamine. In the case of hexamethylenediamine, the dimer product would be bis(hexamethylene)triamine (BHMT). As taught in U.S. Pat. No. 4,739,035 and U.S. Pat. No. 5,627,257 the diamine dimerization is promoted by the high temperatures and pressures required by the processes. These as formed diamine dimers or triamines serve as crosslinkers for the linear polyamide chain. The resultant gel content has been found to lead to significant product quality deterioration.\nThis shortcoming was partially overcome through a two step process disclosed in U.S. Pat. No. 4,739,035. The first step involved the hydrolysis of adiponitrile with water, catalysts and from 0 to 10 weight percent of the total hexamethylenediamine needed for the reaction. The second step consisted of adding the remainder of the hexamethylenediamine followed by polymerization. Such a two step process provided triamine levels in the 560 to 1,300 ppm levels versus the 1,420 to 1,610 ppm levels found through the teachings of other background art.\nThis shortcoming was more fully overcome through a further two step process disclosed in U.S. Pat. No. 5,627,257. The first step consisted of the nearly full hydrolysis of the adiponitrile to adipic acid utilizing a catalyst and a cocatalyst. The cocatalyst was described as a saturated aliphatic or aromatic dicarboxylic acid. The second step consisted of the addition of at least an equimolar amount of diamine followed by polymerization. Through use of this process, they were able to achieve triamine levels between 500 and 600 ppm. A shortcoming of this process was the substantial amounts of catalyst utilized. The level of the catalyst exemplified ranged from about 1 to 1.5 weight percent based on the adiponitrile level. For example, the patent teaches the use of 5 grams of orthophosphorous acid and 1.8 grams of calcium hypophosphite catalyst in the hydrolysis of 487 grams of adiponitrile, (see U.S. Pat. No. 5,627,257, Example 1)."}
{"text": "Field of the Invention\nThe present invention relates generally to low-twist-angle nematic liquid crystal cells with optical compensation for wide-viewing-angle applications, and more particularly pertains to the design of liquid crystal (LC) displays (LCDs) and, particularly, with techniques for maximizing the field of view (or viewing angle) of LCDs while maintaining high contrast ratio viewing from near normal incidence and minimizing variance in relative gray levels over a wide range of viewing angles. These goals are achieved by using low-twist-angle (twist angle from about 30.degree. to about 85.degree.) nematic liquid crystal cells with optical compensation films which have a negative birefringence and tilted optical axes away from the planes of the optical compensation films.\nIn the following description of the present invention, TN (twisted nematic) and LTN (low angle twisted nematic) are used to represent twisted nematic liquid crystal displays with a twist angle of 90.degree. or lower than 90.degree., respectively. For LTN, this usually means a twist angle from about 85.degree. to about 30.degree.. A TN display was invented by Schadt and Helfrich in 1971 (M. Schadt and W. Helfrich, Appln. Phys. Lett. V. 18, 127 (1971)). A liquid crystal (LC) cell, in general, consists of two substrates forming a cavity between them to contain the nematic LC mixture. Between each substrate and the LC medium, there exist a conductive electrode coated on the substrate and an LC alignment layer in direct contact with the LC medium to align the adjacent LC directors into one direction. For a TN LC cell, the direction of LC directors adjacent to one of the substrates is orthogonal (or 90.degree.) to the direction of LC directors adjacent to the other substrate so that the LC directors in the cell twist 90.degree. from one substrate to the other.\nBecause of their maturity in manufacture and their sufficiency in performance, TN liquid crystal (LC) displays have been widely used in commercial thin-film-transistor (TFT) driven flat panel liquid crystal displays (LCDs). The strong viewing-angle dependencies of the contrast ratio, the brightness, and the grayscale of the TFT-driven TN displays have been recognized as major weaknesses for these displays. To illustrate the viewing-angle problem of the TN LC cells, a TN orientation is defined with different viewing directions.\nFIG. 1 shows a TN cell 10 with two substrates 12 and 14. A rubbed polyimide film (not shown in FIG. 1) is usually used to align the LC directors. The rubbing directions of the polyimide films on substrates 12 and 14 are shown as dashed arrow 16 and solid arrow 18, respectively.\nFor display applications, the TN cell is placed between two polarizers 20, 22 with the transmitting axes of the polarizers being either parallel or perpendicular to the adjacent LC directors. If the transmitting axes 24, 26 of the two polarizers sandwiching the TN cell are crossed to each other, the TN display is operated in a normally-white case where the quiescent state of the TN cell is the bright state of the display. On the other hand, if the axes 24, 26 are parallel to each other, the TN display is operated in a normally-black case where the quiescent state of the TN cell represents the dark state. For the normally-white case, there are two optical Eigen modes, the ordinary-ray (o-) and the extra-ordinary-ray (e-) modes, in which the optical field propagates either parallel or perpendicular to the nematic LC directors in the TN cell, respectively. Such e- and o- modes are illustrated in FIG. 1 where the transmitting axes of the polarizers 20 and 22 are shown.\nWith the configuration of the TN display shown in FIG. 1, by facing the display, four viewing zones can be defined, an upper viewing zone for viewing from the 12 o'clock direction, a lower viewing zone for viewing from the 6 o'clock direction, a left viewing zone for viewing from the 9 o'clock direction, and a right viewing zone for viewing from the 3 o'clock direction. The sign of angle for the upper and right viewing zones are positive while the sign of angle for the lower and left viewing zones are negative. Traditionally, the o-mode has been used for bi-level displays. Recently, Takano, et al have carried out a detailed comparison between the o-and e-modes of NW, first-minimum TN cells for analog-grayscale full color displays (H. Takano, M. Ikezaki, and S. Suzuki, \"Threshold Voltage Biased E-mode TN LCD-Optimum Optical Design for Grayscale Applications,\" the IV International Topical Meeting on Optics of Liquid Crystals, Oct. 7-11, 1991, Cocoa Beach, Fla.). They paid particular attention to optimizing the angular region that preserves a proper grayscale order (no grayscale reversal), i.e., minimizing the angular region of grayscale reversals for ratios of eight gray levels. They concluded that the e-mode with a near threshold-voltage bias is superior to the o-mode for analog-grayscale applications. For the rest of the description of the present invention, the e-mode is used as an example. The results are applicable to the o-mode as well.\nTo illustrate the viewing-angle problem of TN for analog-grayscale displays, FIG. 2 shows transmittance as a function of applied voltage for a typical TN cell when the TN cell is being viewed from five different directions. Curves 1, 2, 3, 4 and 5 in FIG. 2 correspond respectively to viewing from normal incidence, 40.degree. from the left viewing zone, 50.degree. from the right viewing zone, 30.degree. from the lower viewing zone, and 30.degree. from the upper viewing zone, where the angles in degrees are defined as the angles of viewing direction with respect to a normal to the display panel. FIG. 2 illustrates that, at a given voltage applied to the TN cell, the brightness (or the contrast ratio) of the display appears different from the above-mentioned five different viewing directions.\nTo further quantify the viewing-angle problems of a typical TN cell, eight different voltage levels were applied to the TN cell to achieve eight approximately equally-spaced gray levels starting from the brightness to the darkest states of the display. The change of these eight levels as a function of viewing angles in the horizontal and vertical viewing directions are shown in FIGS. 3 and 4, respectively, for a typical TN cell. As shown in FIG. 3, at a horizontal viewing angle of either +40.degree. or -40.degree., the transmittance of the gray level 8 (darkest level near normal incidence) is higher than that of the gray level 7. In this case, we have contrast or gray-level reversal between gray level 8 (g8) and gray level 7 (g7) for these viewing directions. The display will appear annoying if a grayscale reversal occurs between any two gray levels from level 1 to level 8.\nFIG. 5 shows iso-contrast curves as a function of viewing angle for a typical TN display. It can be seen that the contrast ratio decreases when the viewing angle deviates further from normal incidence. The TN cell usually has the best contrast ratio near normal incidence. FIG. 5 also shows that, outside the thick solid curves, image (or grayscale) reversal occurs so that the display appears annoying when viewed from these viewing zones having image reversals.\nThe narrow viewing-angle characteristics of a TN cell are caused by a change of retardation of the TN cell when the viewing direction is changed. To a first degree of approximation, the retardation of a nematic LC cell at normal incidence is proportional to d.DELTA.n, where d is the cell gap and .DELTA.n is the birefringence of the LC medium. When the LC cell is viewed from an oblique direction, the retardation becomes larger because the effective cell-gap becomes larger at this oblique viewing direction. Therefore, it is expected that the viewing-angle characteristics of a TN cell can be improved if one uses a low d.DELTA.n for the TN cell. However, as the d.DELTA.n of the TN cell is lower than approximately 0.48 .mu.m, the brightness of the display is reduced. In order to maintain high brightness, we can reduce the twist angle of the LC cell from 90.degree. as the d.DELTA.n is reduced. Therefore, LTN cells can have wider viewing angles than TN cells because the LTN cells have lower values of d.DELTA.n than TN cells for about the same brightness. The major problem associated with LTN cells is that the contrast ratio around normal incidence decreases as the twist angle of the LTN cell decreases. Therefore, a higher operating voltage is required for LTN cells to achieve similar contrast ratios at normal incidence as the regular TN cells. Higher operating voltage implies higher cost and larger power consumption. In other words, LTN cells may have about the same brightness as TN cells and have wider viewing angles than TN cells but the best viewing zone for LTN cells is no longer along the direction of normal incidence. Since the displays are most frequently viewed along the normal incidence, it is important to shift the best viewing zone of LTN cells into near normal incidence without raising the operating voltage.\nOne prior art approach by Hirakata, et al. compensates a 70.degree.-twist LTN cell with uni-axial optical compensation films having positive birefringence (J. I. Hirakata, H. Abe, I. Hiyama, K. Kondo, SID 95 DIGEST, (1995) p. 563.). FIG. 6 illustrates this prior art approach. Imagine the paper containing FIG. 6 being the display. We view the display just as we view FIG. 6, and the incident light impinges from the back side of the display with the viewer situated at the front side. The transmitting axes of the polarizer at the back side and the analyzer at the front side are parallel to the x-axis and y-axis, respectively, as shown in FIG. 6. The LC directors align along the n.sub.1 direction adjacent to the back substrate. The LC directors align along the n.sub.2 direction adjacent to the front substrate. The LC directors in the LTN cell twist from the n.sub.1 direction toward the n.sub.2 direction with a total twist angle. The angle, .alpha., between n.sub.1 and the x-axis is approximately equal to the angle (also .alpha.) between n.sub.2 and the y-axis. From FIG. 6, it can be seen that 2.alpha.+ is approximately equal to 90.degree.. Between the polarizer and the LC cell, there is one sheet of optical compensation film. The optical compensation film is uni-axial with its optical axis along the direction c.sub.1. The direction of c.sub.1 is perpendicular to the direction of n.sub.1. There is another sheet of similar optical compensation film between the LC cell and the analyzer. The optical axis of this second compensation film is along the c.sub.2 direction. The direction of c.sub.2 is also perpendicular to the direction of n.sub.2. Both c.sub.1 and c.sub.2 are in the plane of the compensation film. The d.DELTA.n of the 70.degree.-twist cell is about 0.38 .mu.m, and each compensation film has a positive birefringence and a retardation about 0.023 .mu.m. Retardation films with such a small retardation of about 0.023 .mu.m are expensive because it is difficult to maintain a uniform small retardation within about .+-.3% across the whole display area in order to achieve a display with good uniformity."}
{"text": "1. Field of the Invention\nThe present invention generally relates to generating satellite tracking information for earth orbiting satellites. More specifically, the invention relates to a method and apparatus for generating satellite tracking information in a first format (e.g., a compact ephemeris model) through a network or communications link, then representing the satellite tracking information in a second format (e.g., a standard ephemeris model) at a receiver.\n2. Description of the Related Art\nA positioning receiver for the Global Positioning System (GPS) uses measurements from several satellites to compute a position. The process of acquiring the GPS radio signal is enhanced in speed and sensitivity if the GPS receiver has prior access to a model of the satellite orbit and clock. This model is broadcast by the GPS satellites and is known as ephemeris or ephemeris information. Each satellite broadcasts its own ephemeris once every 30 seconds. Once the GPS radio signal has been acquired, the process of computing position requires the use of the ephemeris information.\nThe broadcast ephemeris information is encoded in a 900 bit message within the GPS satellite signal. It is transmitted at a rate of 50 bits per second, taking 18 seconds in all for a complete ephemeris transmission. The broadcast ephemeris information is typically valid for 2 to 4 hours into the future (from the time of broadcast). Before the end of the period of validity the GPS receiver must obtain a fresh broadcast ephemeris to continue operating correctly and produce an accurate position. It is always slow (no faster than 18 seconds), frequently difficult, and sometimes impossible (in environments with very low signal strengths), for a GPS receiver to download an ephemeris from a satellite. For these reasons it has long been known that it is advantageous to send the ephemeris to a GPS receiver by some other means in lieu of awaiting the transmission from the satellite. U.S. Pat. No. 4,445,118, issued Apr. 24, 1984, describes a technique that collects satellite orbit information at a GPS reference station, and transmits the information to the remote GPS receiver via a wireless transmission. This technique of providing the ephemeris, or equivalent data, to a GPS receiver has become known as xe2x80x9cAssisted-GPSxe2x80x9d. Since the source of ephemeris in Assisted-GPS is the satellite signal, the ephemeris information remains valid for only a few hours. As such, the remote GPS receiver must periodically connect to a source of ephemeris information whether that information is received directly from the satellite or from a wireless transmission. Without such a periodic update, the remote GPS receiver will not accurately determine position.\nFurthermore, the Assisted-GPS systems typically retransmit the entire ephemeris message to the remote receiver. In many instances, bandwidth or packet size for the transmission of this message is not readily available.\nTherefore, there is a need for a method and apparatus for providing satellite trajectory and clock information to a remote receiver in a compact form.\nThe present invention is a method and apparatus for generating satellite tracking data (STD), then transmitting the data to a remote receiver in a compact form. The STD is derived by receiving at one or more satellite tracking stations the signals from at least one satellite and determining satellite tracking information (STI) through signal processing or by extracting the ephemeris message from the received signals. STI contains present satellite orbit trajectory data and satellite clock information.\nThe STD is reformatted into a compact format and provided to a remote satellite signal receiver via a network or communications system. The receiver converts the compact format into a standard format and uses the STD to compute the position of the receiver. The satellite system may include the global positioning system (GPS), GLONASS, GALILEO, or other satellite systems that may use STD to enhance the performance of the receiver."}
{"text": "Various types of equipment exist for semiconductor processing such as plasma etching, ion implantation, sputtering, rapid thermal processing (RTP), photolithography, chemical vapor deposition (CVD), and flat panel display fabrication processes wherein etching, resist stripping, passivation, deposition, and the like, are carried out. For example, a vacuum processing chamber may be used for etching and chemical vapor deposition of materials on substrates by supplying an etching or deposition gas to the vacuum chamber and by application of radio frequency (RF) energy to the gas. Electromagnetic coupling of RF energy into the source region of a vacuum chamber is conventionally employed to generate and maintain a high electron density plasma having a low particle energy. Generally, plasmas may be produced from a low-pressure process gas by inducing an electron flow which ionizes individual gas molecules through the transfer of kinetic energy through individual electron-gas molecule collisions. Most commonly, the electrons are accelerated in an electric field, typically a radiofrequency electric field produced between a pair of opposed electrodes which are oriented parallel to the wafer.\nPlasma generation is used in a variety of such semiconductor fabrication processes. Plasma generating equipment includes parallel plate reactors such as the type disclosed in commonly owned U.S. Pat. No. 4,340,462, electron cyclotron resonance (ECR) systems such as the type disclosed in commonly owned U.S. Pat. No. 5,200,232, and inductively coupled plasma (ICP) or transformer coupled plasma systems such as the type disclosed in commonly owned U.S. Pat. No. 4,948,458.\nDue to the tremendous growth in integrated circuit production, the use of vacuum processing chambers has increased dramatically in recent years. The use of vacuum processing chambers may seriously affect the environment, however, because perfluorocarbons (PFCs) are widely used in plasma etch and plasma-enhanced CVD equipment. PFCs are highly stable compounds which makes them well suited for plasma processing. However, PFCs also significantly contribute to global warming and are not destroyed by scrubbers or other conventional emission control equipment used in vacuum processing chambers. Although there are many gasses which cause global warming, PFCs and hydrofluorocarbons (HFCs), both of which are referred to hereinafter as “fluorocarbons”, also used in plasma processing, have particularly high global warming potentials (GWPs). For example, CF4, C3F8, SF6, NF3, and C2HF5 all have GWPs of over 3000, and C2F6, SF6, and CHF3 have GWPs of over 12,000. By contrast, carbon dioxide, a well-known greenhouse gas, has a GWP of 1. In addition, because of their stability, PFCs have a very long lifetime. For example, the lifetimes of CF4 and C2F6 are 50,000 and 10,000 years, respectively. Thus, collectively, these process gasses can have a significant impact on the environment.\nTo reduce the impact of PFCs on the environment, several conventional methods for abating PFCs from vacuum processing chambers have been proposed, including process optimization, chemical alternatives, and destruction/decomposition. Process optimization involves the refinement of system parameters to achieve the desired process while using the minimum amount of PFCs. Process optimization is desirable because it reduces chemical costs and emissions and may increase throughput and prolong the life of internal components of the reactor. However, process optimization does not provide a complete solution since it does not involve the abatement of PFCs once they are used in the system. Thus, although the amount of PFCs used is reduced by process optimization, the PFCs that are used are ultimately emitted into the environment.\nChemical alternatives to using PFCs are desirable because they eliminate the problem of PFC emissions entirely. However, to date, research is still underway to uncover more effective and environmentally sound chemical alternatives.\nThere are two basic categories for conventional destruction/decomposition techniques. The first category involves abatement performed on the atmospheric side of the system, either at each tool or on a large scale for multiple tools, after the gasses have passed through the pumping system. On the atmospheric side, there are several possibilities, including water scrubbers, resin beds, furnaces, flame-based burn boxes, and plasma torches. All of these except burn boxes and torches are ineffective against many of the highly stable PFC compounds.\nBurn boxes have been shown to be inefficient abatement devices, unless very large amounts of reactant gases such as hydrogen, methane, or natural gas are flowed through the burn box. This makes these devices very expensive to operate, and environmentally unfriendly.\nPlasma torches on the atmospheric side could be more effective; however, very large and expensive, not to mention dangerous, torch facilities would be required to abate what would be a small concentration of PFCs in a very high flow of tool effluent. This inefficiency is exacerbated by the addition of large amounts of nitrogen, used as a diluent in vacuum pumps and as a purge gas in many tool operations.\nThe second general category of destruction/decomposition techniques involves abatement performed under a vacuum upstream of the pumping system. Plasma destruction is one method, for example, in which a device is employed to treat the exhaust from the tool upstream of the pump before nitrogen purge dilution has taken place. In plasma destruction, energy is applied to the reaction gasses to create a plasma in which the gasses are ionized. The PFCs become unstable at the high energy state, and are consequently broken down into smaller molecules which are less detrimental to the environment.\nExamples of devices for plasma destruction include the ETC Dryscrub and the Eastern Digital Post-Reaction Chamber (PRC). The ETC Dryscrub is a flat spiral chamber. The gases come in at the outer end of the spiral, circle around, and eventually exit through the center of the spiral. The spiral is an RF electrode operated at 100 kHz. The purpose of the reactor is to dissociate exhaust gases coming from a CVD system so that any remaining solid-producing gases are reacted onto the walls of the spiral. Tests performed on the ETC Dryscrub reactor, however, resulted in a relatively ineffective abatement of C2F6, the test gas. In addition, the main reaction product was another greenhouse PFC gas, CF4, and most of the secondary products were also greenhouse gasses. The Eastern Digital PRC is essentially the same as the ETC Dryscrub, and yields similar results. These devices are both inefficient at plasma abatement and have low plasma densities and dissociation rates.\nAnother known method for plasma destruction, in which CF4, C2F6, and SF6 may be abated, involves the use of a microwave plasma reactor. Microwave sources, however, are expensive and complex. They also have small skin depths, so they tend have an axial region where there is no plasma, through which unabated gases escape. This can be compensated by inserting a “plug”, but the plug then reduces the fluid conductance of the device, which in turn adversely affects the performance of the pumping system."}
{"text": "As the development of nanoscale mechanical, electrical, chemical and biological devices and systems increases, new processes and materials are needed to fabricate nanoscale devices and components. Conventional optical lithographic processing methods are not able to accommodate fabrication of structures and features much below a 100 nm level. The use of self-assembling diblock copolymers presents another route to patterning at nanometer dimensions. Diblock copolymer films spontaneously assemble into periodic structures by microphase separation of the constituent polymer blocks after annealing, for example, by thermal annealing above the glass transition temperature of the polymer or by solvent annealing, forming ordered domains at nanometer-scale dimensions. Following self-assembly, one block of the copolymer can be selectively removed and the remaining patterned film used as an etch mask for patterning nanosized features into the underlying substrate. Since the domain sizes and periods (Lo) involved in this method are determined by the chain length of a block copolymer (MW), resolution can exceed other techniques such as conventional photolithography, while the cost of the technique is far less than electron beam lithography or EUV photolithography, which have comparable resolution.\nThe film morphology, including the size and shape of the microphase-separated domains, can be controlled by the molecular weight and volume fraction of the AB blocks of a diblock copolymer to produce lamellar, cylindrical, or spherical morphologies, among others. For example, for volume fractions at ratios greater than about 80:20 of the two blocks (AB) of a diblock polymer, a block copolymer film will microphase separate and self-assemble into a periodic spherical domains with spheres of polymer B surrounded by a matrix of polymer A. For ratios of the two blocks between about 60:40 and 80:20, the diblock copolymer assembles into a periodic hexagonal close-packed or honeycomb array of cylinders of polymer B within a matrix of polymer A. For ratios between about 50:50 and 60:40, lamellar domains or alternating stripes of the blocks are formed. Domain size typically ranges from 5 nm to 50 nm.\nDiblock copolymer thin films of cylindrical and lamellar phases may both form striped phases relative to an interface. For cylindrical phase films, a striped pattern results from parallel cylinder orientation, while for lamellar phase films, a striped pattern results from perpendicular domain orientation. From a top down view, perpendicular-oriented lamellae and parallel-oriented cylinders appear similar, e.g., as parallel lines.\nGraphoepitaxy techniques using defined topography such as trench edges have been used in an attempt to orient and order copolymer domains and control registration and alignment of the self-assembled blocks to form a desired pattern. By comparison, thin films prepared on a flat substrate and annealed do not have any preferential orientation of domains and will assemble into a disordered fingerprint-like structure despite extensive annealing. Although registered and ordered arrays of cylinders have been produced within trenches, the fabrication of an ordered array of repeat structures outside of the confines of the trenches and over large areas has not been realized in a manufacturable process.\nIt would be useful to provide a method of fabricating films of linear arrays of ordered nanostructures that overcome these problems."}
{"text": "Digital pens and related technologies have been gaining popularity due to their convenience and versatility. One of the advantages provided by digital pens is the ability to store handwritten information in the pen itself, eliminating the need to collect, transport, and/or store individual pieces of paper upon which such information may be written or to immediately download the stored information to a computing device.\nHowever, this advantage also leads to one of the primary disadvantages associated with digital pens, namely, the potential for confidential or otherwise sensitive information to be compromised by the loss or unauthorized accessing of the digital pen. Thus, information stored in a digital pen may be unnecessarily exposed to compromise. This may violate not only good judgment and individual or business guidelines, but may also expose the user of the digital pen to liability for violation of local, state, or federal law (e.g., Health Insurance Portability and Accountability Act of 1996 (HIPAA), Sarbanes-Oxley Act of 2002, etc.).\nAccordingly, there exists a need in the art to overcome the deficiencies and limitations described hereinabove."}
{"text": "1. Field of the Invention\nThe invention disclosed herein relates to evaluating reservoirs in earth formations, and, in particular, to modeling flow properties of the reservoirs.\n2. Description of the Related Art\nHydrocarbons are typically recovered by having the hydrocarbons flow out of reservoirs in formation rock and into a well that penetrates the formation rock. The hydrocarbons are then extracted from the well at the surface of the earth.\nIn order to efficiently use production resources, a commercial reservoir simulator may be used to model flow in fractured media. Commercial reservoir simulators convert permeable fracture networks into equivalent porous media for simulation of the fractured media. Unfortunately, these commercial reservoir simulators lose accuracy due to their use of approximations. It would be well received in the hydrocarbon recovery industry if models of the hydrocarbon containing reservoirs could be improved to increase their accuracy."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the surface mounting of ceramic leadless chip carriers to printed wiring boards. More particularly, the present invention relates to methods and apparatus for increasing the reliability of the solder joint integrity between the printed wiring board and the chip carrier.\n2. Description of Related Art\nThe surface mounting of electronic devices, such as ceramic leadless chip carriers, has become popular due to the ability to save weight and volume in electronic equipment when such surface-mount technology is used. In addition, the need for increased circuit density makes surface mounting of chip carriers even more desirable.\nTypically, the ceramic leadless chip carrier is mounted to the printed wiring board surface by way of solder fillets or joints. It is important that the integrity and reliability of the solder joint be maintained for long periods of time and be able to withstand numerous temperature and power cycles. In addition, the solder joint must be able to withstand shock and vibration which may occur in many electronic systems such as those used in avionics.\nThe increasing popularity of surface mounting of ceramic leadless chip carriers has brought with it concerns over the reliability and integrity of the solder joints between the chip carriers and the printed wiring board. Cracking and other disturbances to the joint integrity have been noted. Joint cracking is believed to be due mainly to differences in the thermal coefficients of expansion between the printed wiring board and the chip carriers. Various different approaches have been taken to reduce the possibility of solder joint cracking. These approaches have included making the printed wiring board out of materials having thermal coefficients of expansion which are close to the thermal coefficient of expansion for the ceramic leadless chip carriers. In addition, attempts have been made to replace the solder joint with a joint which is less likely to crack, such as electrically conductive epoxy compounds. Although the various approaches used to date have experienced different degrees of success, there still is a continuing need to provide a method for surface mounting ceramic leadless chip carriers to printed wiring boards wherein the possibility of solder joint cracking is eliminated."}
{"text": "The present invention is directed to optical transmission lines of the type that normally couple one or more utilization devices into a signal loop and more particularly, to a by-pass relay for connecting and disconnecting a utilization device from an optical transmission line while still maintaining the optical continuity of the line.\nOptical transmission lines have come into recent use because of their ability to couple signals from one point to another with a minimum amount of interference.\nIn electrical systems, interference is generally caused by electromagnetic radiaton, cross-coupling of signals and/or noise. These interfering effects, of course, diminish the quality of the received signal. In digital systems, for example, computer systems, the requirement for coupling a central processing unit to one or more peripheral units located in different cabinets some distance from the central processing unit creates a requirement for coupling signals with interference eliminated totally, or minimized to the smallest of levels, otherwise computing errors will result.\nIn such systems, it is necessary that the peripheral units be capable of being logically connected or disconnected from the transmission line without disrupting the operation of the loop as it affects other peripheral devices.\nFrom the foregoing, it appears natural that a transmission line that would have the best noise immunity would be implemented using fiber optics. Such being the case, the present invention directs itself to the problem that arises as regards the connecting and the disconnecting of the peripheral device into the optical transmission line. Optical transmission lines may consist of a single optical fiber or a bundle of optical fibers. The interconnections between remotely disposed stations is established utilizing schemes that are analogous to electrical systems. These stations may be interconnected by a loop or a line data bus which interconnection reduces the number of optical signal transmission lines that have to be used from, for example, the number that would be required if each individual station and/or peripheral device were connected directly to the central station or computer processing unit CPU by its own transmission path."}
{"text": "The background is presented herein only by way of example to multichip modules (MCMs), which are multimaterial, multilayered devices that are becoming one of the electronics packaging industry's most preferred components for a variety of aerospace, computer, military, and telecommunications applications. MCMs are replacing or reducing the complexity of printed circuit boards, thus enhancing product efficiency and reliability.\nMCMs and other multimaterial, multilayered electronic devices for packaging single chips such as ball grid arrays (BGA), pin grid arrays (PGA), etc; circuit boards; and hybrid and semiconductor microcircuits typically include separate component layers of metal and an organic dielectric and/or reinforcement materials. The standard metal component layer(s) may contain aluminum, copper, gold, molybdenum, nickel, palladium, platinum, silver, titanium, or tungsten, or combinations thereof. These layers typically have a depth or thickness of about 9 to 36 .mu.m (where 7.8.times.10.sup.-3 kg of metal equals a thickness of about 9 .mu.m), but may be thinner or as large as 72 .mu.m. A standard organic dielectric layer may include bismaleimide triazine (BT), cardboard, cyanate esters, epoxies, phenolics, polyimides, or polytetrafluorethylene (PTFE). These layers typically have a depth of about 50 to 400 .mu.m. A standard reinforcement component \"layer\" may include fiber matts or dispersed particles of aramid fibers, ceramics, glass, or Kevlar.TM. woven or dispersed into the organic dielectric layer to reinforce it. These reinforcements typically have a diameter or thickness of about 1 to 10 .mu.m. Stacks having several layers of metal, dielectric, and reinforcement material may be larger than 2 mm.\nMCMs present, however, new manufacturing obstacles because they require finer lines and smaller high aspect ratio interconnections or vias and use a variety of new materials. The term \"vias\" is used herein generally to refer to complete through-holes or incomplete holes called \"blind vias.\" A high aspect ratio implies a depth-to-width comparison that is typically from about one to five. After the vias are formed, they are typically plated with a metal to form interconnections between metal circuits in the multilayered electronic devices. The ability to increase the packaging density of MCM-L, MCM-L/D, PCMCIA, and laminate BGA devices to fulfill present and future technology requirements is now limited by existing processes for forming interconnections between layers. More detail concerning the industry trend and its obstacles may be found in \"The Interconnect Challenge: Filling Small, High Aspect Ratio Contact Holes,\" Semiconductor International, Vol. 17, No. 9, Aug. 1994, at 57-64.\nThe laminate circuit board industry has explored a variety of multistep processes for forming interconnections. The most common interconnection formation process includes via formation, precleaning, surface activation, preplating, plating, and homogenization or planarization. Via formation is generally achieved by mechanical drilling or punching, but has also been accomplished to a limited extent by excimer and CO.sub.2 laser ablation. The precleaning or desmearing process typically entails a chemical process employing an acid or permanganate solution.\nThe surface activation process generally entails seeding the internal wall surfaces with palladium to form a preliminary conductive layer or a metal adhesive layer. The palladium can be directly deposited as ions from an alkaline solution or in the form of an organometallic coating, the organics of which are subsequently chemically removed. These processes seem to work in the range of about 150 to 250 .mu.m. Other techniques, such as carbonizing the internal wall surfaces of the via to make them sufficiently conductive for certain preplating processes, are sometimes employed.\nElectroless chemical plating is the most common method for preplating the internal wall surfaces of vias. Electroless chemical plating employs a liquid solution of electroless copper to chemically react with and deposit on the internal wall surface of the via a few nanometers thickness of metal. This type of copper plating reduces the resistance across the outer layers of the multilayered electronic device to about 10 ohms (for thousands of vias). The chemical reaction generates a gas byproduct (hydrogen and oxygen) having a volume that is approximately six times the volume of a via. The electronic device is typically agitated in solution to exchange gas and reactant inside the via. At via diameters smaller than about 250 .mu.m, the gas and reactant exchange becomes economically impractical and generally infeasible, resulting in poor coverage of the via internal wall surface.\nA typical electroless copper process line includes approximately 30 tanks, each containing over 100 gallons of chemical solutions, that are sequentially connected together by a conveyor handling system. The process includes precleaning, surface activation, and electroless copper deposition steps. The entire process line is typically 60-100 feet long, and the equipment costs approximately $250,000. The various chemical solutions may cost up to about $300/gallon.\nElectrochemical plating is the industry-preferred method for plating vias. The electrochemical plating process places a charge across the outer layers of the via while the electronic device is submersed in a plating bath of a copper solution. The copper ions are galvanically attracted to and cling to the via internal wall surfaces. This process is limited by its speed and the access that the electrolyte solution has to the internal wall surfaces. Thus, the vias are not totally filled. A variety of other methods for creating electrical interconnects adopted from the hybrid circuit and semiconductor industries, some of which are described below, have also been used for preplating and then followed by electrochemical plating.\nPlating has also been accomplished by application of a variety of organometallic solid- or liquid-phase compositions. The compositions, such as conductive pastes, are either squeezed into the vias or the vias are submersed in the baths of the compositions. The organics are then chemically or thermally removed or displaced with more or other metals. This process generally creates porous interconnections.\nPhysical vapor deposition (PVD), a process primarily used for creating thin layers on horizontal surfaces, has also been adapted for plating vias. A PVD process converts the plating material into the vapor phase, transports the vapor across a region of reduced pressure, and condenses the material onto the horizontal surface to form a thin film. The two common types of PVD, evaporative PVD and sputtering, employ the addition of heat or the physical dislodgement of surface atoms by momentum transfer, respectively, to convert the material to the vapor phase. The processes necessitate a protective mask to safeguard the top surface of the device prior to via plating or an etching step to remove undesirable metal after via plating. These processes also cause the formation of \"keyhole voids.\" These voids occur because vias typically have aspect ratios greater than one, so the top via rim accumulates metal and becomes sealed before the entire via volume is filled. In some cases, especially whenever the aspect ratio is high, the vias may inadvertently be sealed prior to the sufficient coating of the entire internal wall surfaces.\nThe chemical vapor deposition (CVD) process, on the other hand, reacts chemical vapors of desired concentrations with the heated surface of the device to form a thin film. Like the other vapor deposition processes, CVD necessitates device surface protection when the process is utilized as a via plating technique.\nPostplating homogenation or planarization are typically employed to diffuse metals into porously filled vias or to reduce the voids or diminish the lips formed at via rims. These processes are generally thermal.\nSome researchers, particularly in the semiconductor industry, have investigated using laser technology to assist in some of the steps for forming a conductive interconnect. For example, lasers have been used to modify the surface in preparation for a metallization step, a process similar to that employed in copy machine technology. The surface is charred with, for example, a CO.sub.2 laser and then exposed to some standard thin film technique that preferentially responds to the laser- activated sites. Lasers have also been used to enhance chemical vapor deposition where the laser acts to heat the via site and prepare the area for the chemical diffusion of the metals in the vapor (pyrolytic decomposition). This method is described in U.S. Pat. No. 5,060,595 of Ziv et al., but this method works only with certain Si and Si oxide combinations that exhibit high absorption at 532 nm (frequency-doubled YAG).\nA variety of laser metal deposition techniques has also been explored for use in direct imaging of circuits as described by Liu in \"Laser Metal Deposition for High-Density Interconnect,\" Optics & Photonics News, June 1992, at 10-14. Such laser-activated chemical deposition processes are generally characterized as either photolytic or pyrolitic. These processes are typically followed by chemical electroplating techniques.\nLaser photolytic processes employ laser light to directly cause chemical changes in organo-metallic solutions such that the metals dissociate from the organics and cling to the internal wall surfaces of the vias. These photolytic processes generally result, however, in plating speeds that are too slow for commercial use.\nLaser pyrolytic processes employ heat from a laser source to break apart organo-metallic molecules applied as thin films to electronic devices. The heat releases the organic molecules and leaves the metals. A variety of metals have been deposited in this manner.\nAnother method uses an XeCl excimer laser at 308 nm to planarize a top aluminum layer causing melted aluminum to flow into vias creating interconnections as described by Bachmann in \"Physical Concepts of Materials for Novel Optoelectronic Device Applications,\" SPIE, Vol. 1361, 1990, at 500-11.\nExcimer lasers have also been employed to enhance PVD sputtering processes at very low power and repetition rates to produce thin films of several materials. Further details concerning this technique can be found in \"Pulsed Laser Deposition of Thin Metallic Alloys,\" Appl. Phys. Lett., Vol. 62, No. 19, May 10, 1993, by Krebs and Bremert.\nSkilled persons will appreciate that the cleaning, preplating, and planarizing techniques described above are all relatively slow batch processes, requiring expensive baths, masks, or postprocessing. In addition to being expensive and largely inefficient, most of the plating techniques require both a cleansing as well as a preplating process. The industry needs, therefore, faster and more efficient via plating and/or preplating processes that minimize the tooling as well as the steps for plating vias comprising a wide variety of sizes and internal wall materials."}
{"text": "Starches from grain, cereals, and tubers, e.g., cornstarch, are widely used in the industrial manufacture of products such as sugar syrups and biofuels. For example, high fructose corn syrup (HFCS) is a processed form of corn glucose syrup having high fructose content and a sweetness comparable to sucrose, making HFCS useful as a sugar substitute in soft drinks and other processed foods. HFCS production currently represents a billion dollar industry. Similarly, the production of ethanol from starches is a rapidly expanding industry.\nSyrups and biofuels can be produced from starch by an enzymatic process that catalyzes the breakdown of starch into glucose. This enzymatic process typically involves a sequence of enzyme-catalyzed reactions:\n(1) Liquefaction: Alpha-amylases (EC 3.2.1.1) first catalyze the degradation of a starch suspension, which may contain 30-40% w/w dry solids (ds), to maltodextrans.\nAlpha-amylases are endohydrolases that catalyze the random cleavage of internal α-1,4-D-glucosidic bonds. Because liquefaction typically is conducted at high temperatures, e.g., 90-100° C., thermostable alpha-amylases, such as alpha-amylases from Bacillus sp., are preferred for this step. Alpha-amylases currently used for this step, e.g., alpha-amylases from B. licheniformis, B. amyloliquefaciens, and Geobacillus stearothermophilus (AmyS), do not produce significant amounts of glucose. Instead, the resulting liquefact has a low dextrose equivalent (DE), containing maltose and sugars with high degrees of polymerization (DPn).\n(2) Saccharification: Glucoamylases catalyze the hydrolysis of alpha-1,4-glucosidic linkages of maltodextrins formed after liquefaction from non-reducing ends, releasing D-glucose. Saccharification produces high glucose syrup. Debranching enzymes, such as pullulanases, can aid saccharification.\n(3) Further Processing: A branch point in the process occurs after the production of a glucose-rich syrup. If the final desired product is a biofuel, yeast can ferment the glucose-rich syrup to ethanol. On the other hand, if the final desired product is a fructose-rich syrup, glucose isomerase can catalyze the conversion of the glucose-rich syrup to fructose.\n\nAlpha-amylases are isolated from a wide variety of bacterial, fungal, plant, and animal sources. Many industrially important alpha-amylases are isolated from Bacillus sp., in part because of the generally high capacity of Bacillus to secrete amylases into the growth medium. In addition, Bacillus alpha-amylase variants with altered while more desirable properties are obtained through genetic engineering. Furthermore, there is a need for blends of alpha-amylases, or variants thereof, which can capitalize on the best properties of at least two alpha-amylases of different origins.\nThe Fuelzyme®-LF alpha-amylase (SEQ ID NO: 2)(Verenium Corp.) is an engineered alpha-amylase obtained through DNA shuffling of three parental enzymes. See Richardson et al., J. Biol. Chem. 277: 26501-26507 (2002); U.S. Pat. No. 7,323,336. The advantageous properties of the Fuelzyme®-LF alpha-amylase include: effective viscosity reduction at a lower dose, improved thermostability, and broad pH operating ranges. The use of this alpha-amylase, however, is currently limited to biofuel applications, e.g., ethanol production, because it results in ineffectual glucose syrup that is not suitable for downstream applications such as sweetener applications. Specifically, saccharification of starch liquefact from Fuelzyme®-LF alpha-amylase results in iodine-positive saccharide (IPS), which indicates incomplete starch hydrolysis. Thus, if a way could be found to fully exploit the advantages of the Fuelzyme®-LF alpha-amylase in starch processing, particularly in sweetener applications, by using an optimized blend of alpha-amylases, this would also represent a useful contribution to the art."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of computer systems, and more particularly to a multiprocessor computer system having an input/output bridge coupled to devices external to the multiprocessor computer system.\n2. Description of the Relevant Art\nPersonal computers (PCs) and other types of computer systems have been designed around a shared bus system for accessing memory. One or more processors and one or more input/output (I/O) devices may be coupled to the memory through the shared bus. The I/O devices may be coupled to the shared bus through an I/O bridge which manages the transfer of information between the shared bus and the I/O devices, while processors are typically coupled directly to the shared bus or coupled through a cache hierarchy to the shared bus.\nUnfortunately, shared bus systems suffer from several drawbacks. For example, since there are multiple devices attached to the shared bus, the bus is typically operated at a relatively low frequency. The multiple attachments present a high capacitive load to a device driving a signal on the bus, and the multiple attach points present a relatively complicated transmission line model for high frequencies. Accordingly, the frequency remains low, and the bandwidth available of the shared bus is similarly relatively low. A low bandwidth presents a barrier to attaching additional devices to the shared bus, as performance may be limited by available bandwidth.\nAnother disadvantage of the shared bus system is a lack of scalability to larger numbers of devices. As mentioned above, the amount of bandwidth is fixed and may decrease if additional devices are added. Once the bandwidth requirements of the devices attached to the bus, either directly or indirectly, exceeds the available bandwidth of the bus, devices will frequently be stalled when attempting to access the bus. As a result, overall performance may be decreased.\nOne or more of the above problems may be addressed using a distributed memory system. A computer system employing a distributed memory system includes multiple nodes. Two or more of these nodes are connected to individual memories, respectively, and the nodes are interconnected using any suitable interconnect. For example, each node may be connected to each other node using dedicated lines. Alternatively, each node may connect to a fixed number of other nodes, and transactions between nodes may be routed from a first node to a second node to which the first node is not directly connected via one or more intervening nodes. A memory address space is assigned across the memories in each node.\nNodes may additionally include a processor. The processor typically includes a cache which stores cache blocks of data read from the memories. Furthermore, a node may include one or more caches external to the processors. Since the processors and/or nodes may be storing cache blocks accessed by other nodes, it is desirable to maintain coherency within the nodes.\nThe present invention provides a method and circuit for maintaining order of memory requests initiated by I/O devices coupled directly or indirectly to a multiprocessor computer system via an I/O bridge. In one embodiment, the multiprocessor computer system comprises a plurality of processing nodes and a plurality of memories. Each processing node includes at least one microprocessor coupled to a memory controller which, in turn, is coupled to one of the plurality of memories. Each memory defines a portion of the memory address space for the computer system. Each processing node is uniquely identified by a node number. The I/O bridge is coupled to one or more processing nodes and is configured to generate and transmit a non-coherent memory access transaction. The non-coherent memory access transaction is generated in response to a memory access request generated by an I/O device coupled directly or indirectly to the I/O bridge. The I/O bridge transmits the non-coherent memory access transaction to at least one processing node connected thereto where it is transformed into a coherent memory access transaction for transmission to another processing node within the multiprocessor computer system. Each non-coherent memory transaction includes at least one command packet containing a pipe identification and a memory address affected by the memory access transaction. The pipe identification may include information identifying the origin of the memory request. Generally, non-coherent memory transactions having the same pipe identification must complete in order whereas non-coherent memory transactions having different pipe identifications may complete out of order.\nThe I/O bridge generates first and second non-coherent memory access transactions, the second non-coherent memory access transactions being generated after the first non-coherent memory access transaction. Each of the first and second non-coherent memory access transactions include first and second memory addresses, respectively. Further, each of the first and second non-coherent memory access transactions include first and second pipe identifications, respectively. The first and second memory addresses are mapped to first and second node numbers, respectively, in response to the connected node receiving the first and second non-coherent memory access transactions. Thereafter, the first and second pipe identifications of the received first and second non-coherent memory access transactions, are compared. Additionally, the first and second node numbers are compared if the first and second pipe identifications compare equally. The connected node which receives the first and second non-coherent memory transactions generates and transmits a first coherent memory access transaction corresponding to the first non-coherent memory transaction. This first coherent memory access transaction includes the mapped first node number and the first memory address. The processing node also generates a second coherent memory access transaction corresponding to the second non-coherent memory access transaction. This second coherent memory access transaction includes the second node number and the second memory address. However, in one embodiment, the second coherent memory access transaction is generated only if the first and second pipe identifications do not compare equally or if the first and second node numbers compare equally. In another embodiment, the second coherent memory access translation is generated without delay upon the node circuits receipt of the second non-coherent memory access transaction. However, the second coherent memory access transaction is transmitted from the connected node processing only if the first and second pipe identifications do not compare equally or if the first and second node numbers compare equally."}
{"text": "Cushioning materials have a variety of uses, such as for mattresses, seating surfaces, shoe inserts, packaging, medical devices, etc. Cushioning materials may be formulated and/or configured to reduce peak pressure on a cushioned body, which may increase comfort for humans or animals, and may protect objects from damage. Cushioning materials may be formed of materials that deflect or deform under load, such as polyethylene or polyurethane foams (e.g., convoluted foam), vinyl, rubber, springs, natural or synthetic fibers, fluid-filled flexible containers, etc. Different cushioning materials may have different responses to a given pressure, and some materials may be well suited to different applications. Cushioning materials may be used in combination with one another to achieve selected properties.\nFor example, cushioning materials may include a foam layer topped with a layer of thermoset elastomeric gel, such as a polyurethane gel or a silicone gel. Because polyurethane gels and silicone gels are generally structurally weak and/or sticky, cushioning materials may include film covering such gels, such as a thin thermoplastic polyurethane film. The film may reinforce the strength of the gel, and may prevent other materials from sticking to the gel, since the film generally adheres to the gel but is not itself sticky.\nGels may be used for cushioning and/or temperature management. Gels may provide cushioning because the gels may hydrostatically flow to the shape of a cushioned object and may tend to relieve pressure peaks. Gels may also reduce stresses from shear. Gels may have high thermal mass and/or thermal conductivity, and may therefore be used for heating (such as in hot packs for sore muscles), cooling (such as in cold packs for sprains or for a feeling of coolness when lying on a mattress or pillow), or maintaining a given temperature (such as in a mattress being used in a warm or cool room). For example, gel may be fused to the top of a mattress core, and a film may cover the gel. As another example, gels may be used as the top layer of a gel-on-foam wheelchair cushion.\nA conventional gel layer, with or without a plastic film, may be a barrier to gases (e.g., air, vapors, or other gases). This barrier may cause difficulties such as discomfort, such as when body heat and/or perspiration accumulate between the user's body and the gel layer. Even when a breathable material (such as a cover comprising foam or batting fiber) is disposed between a cushioned object and the gel, gases can only travel laterally through the breathable material. Since gases cannot penetrate the plastic film or the gel, the plastic film or the gel inhibits the flow of the gases away from the cushioned object. When the weight of the cushioned object compresses the breathable material, the lateral gas flow paths may become more constricted. Thus, it would be beneficial to provide a cushioning material that alleviates some of these concerns."}
{"text": "Electronic components for audio and/or audiovisual applications conventionally include multiple, free-standing enclosures that receive power and signals from facility wiring and communicate with other components on wired cables or wireless links. Support for numerous components has conventionally been provided by furniture called an entertainment center. A conventional entertainment center may have open shelving and/or enclosed shelving for supporting and enclosing not only the components but also media used with the components. Such furniture also conventionally provides holes through the back and through the shelving for accommodating the signal cables and power cables associated with the components.\nA conventional entertainment center is spaced away from a facility wall to allow cabling to be tucked behind the cabinetry of the entertainment center because provisions for cabling inside the cabinetry of the entertainment center are inadequate. The space between the entertainment center and the facility wall also supplies ventilation air for the components.\nConventional entertainment centers typically provide movable shelving for accommodating electronic components of different vertical height, however, such centers typically provide a fixed horizontal dimension to accommodate a maximum component width.\nThus, use of a conventional entertainment center is limited by the fixed horizontal width of its design. Users of such conventional entertainment centers seeking, for example, to accommodate a larger home theater display (e.g., a big screen television set, a rear projection system, or a front illuminated screen) have little recourse but to purchase new furniture in the event the larger width display does not fit the fixed horizontal width provided by the existing entertainment center.\nA large market exists for furniture to support objects such as electronic components, e.g., audio system components and/or audio visual system components which may include stereo system components, television system components, home theater system components and the like. New electronic components and other products of various sizes are frequently launched. Consumers owning conventional entertainment centers for accommodating objects having a horizontal width less than a fixed width associated with the conventional entertainment center may be reticent to purchase new entertainment centers for accommodating objects, e.g., electronic components, having a horizontal width greater than the fixed width associated with the conventional entertainment center. Moreover, such consumers may also forego purchasing newer larger electronic components if the existing fixed-width entertainment centers cannot accommodate the horizontal width of the newer larger electronic components. Consequently, without the present invention, both the consumer electronics and furniture industries face economic impairment."}
{"text": "Conventional communications cables typically include a core of twisted pairs of insulated conductors that are enclosed in a protective jacket. To avoid crosstalk with other cables, often referred to as alien crosstalk, conventional cables often include at least one shielding layer disposed around the core of twisted pairs of conductors. A non-conductive barrier layer between the shielding layer and the pairs must also be provided to insulate the core of pairs of the cable.\nSuch cables, however, are often bulky because of the requirement of both a shielding layer and a barrier layer. Therefore, a need exists for a cable that protects against alien crosstalk and that is not bulky."}
{"text": "1. Field of the Invention\nThis invention relates to a process for the stream reforming of methanol, and particularly to a process for producing a reformed gas containing hydrogen gas and carbon dioxide gas in a good efficiency by vapor phase reaction of methanol and water in the presence of catalysts.\n2. Description of the Prior Arts\nHydrogen gas is used in many industries. For instance, it is extensively used in ammonia synthesis, hydrogenation of various types of organic compounds, oil refining and desulfurization, etc. Also, it has a great demand for the metallurgical industry and the semiconductor industry.\nFurthermore, hydrogen gas is expected to be a new energy source accompanying a recent advance in the fuel battery technologies. Thus, the demand for hydrogen gas has increasingly expanded.\nConventional methods for the production of hydrogen gas include steam reforming of hydrocarbons such as liquefied petroleum gas (LPG), liquefied natural gas (LNG) and naphtha. According to the usual method, carbon monoxide and carbon dioxide are removed from a reformed gas containing hydrogen gas, carbon monoxide and carbon dioxide obtained by the above methods so as to produce hydrogen gas. The conventional methods have several drawbacks in that (1) the price of the raw material hydrocarbons continues to rise after the oil shock, and the supply of the raw material hydrocarbons is in unstable conditions, (2) desulfurization of the raw materials is required, (3) a high reaction temperature of 800.degree.-1,000.degree. C. is required for steam reforming, etc. Thus, the conventional methods are suitable for large scale hydrogen gas production, but they are not adequate for middle to small scale hydrogen gas production.\nIn contrast, hydrogen gas production by the steam reforming of methanol has various advantages in that (1) the reaction temperature is relatively low, (2) separation of hydrogen gas from the reformed gas is easy, (3) no desulfurization is required because methanol is the raw material, etc. Also, the method can easily cope with large to small scale plants, since it uses inexpensive and easily transportable methanol as the raw material.\nThe steam reforming reaction of methanol (CH.sub.3 OH+H.sub.2 O.fwdarw.3H.sub.2 +CO.sub.2) itself has long been known, and some patents relating to the reaction are made open. In general, the reaction is carried out by adopting the molar ratio of water to methanol at 1.0 or more. The reaction yields a wet gas (hereinafter referred to as wet gas) containing condensable components such as methanol and water, and a reformed gas. After the reaction, the reformed gas containing hydrogen and carbon dioxide is taken out by cooling of the wet gas. Here, an industrial problem is the treatment of the condensed liquid. Conventionally, the condensed liquid is disposed without treatment as in FIG. 5, or the condensed liquid obtained by vapor-liquid separation treatment is recycled to the reaction system as in FIG. 6 so as to make reuse of the condensed liquid together with the raw material methanol and water. However, the former has a big problem from the view-point of pollution since the condensed liquid contains a considerable amount of organic components such as unreacted methanol and high boiling point components. In the latter, a trace amount of ethanol contained in the raw material methanol accumulates because it is hardly converted by the usual reaction method."}
{"text": "1. Technical Field\nThe present invention relates to pattern density control, and more specifically, to pattern density control using edge printing processes.\n2. Related Art\nIn a conventional sidewall image transfer (SIT) process, memory regions are formed from a memory layer and then used as a blocking mask for etching one or more layers beneath the memory regions so as to transfer the image of the memory regions down to the underlying layer(s). The formation of the memory regions from the memory layer involves a planarization-and-etchback step. However, the planarization-and-etchback step may have variations due to pattern density variation on the wafer. Therefore, there is a need for pattern density control methods in sidewall image transfer (SIT) technologies and other edge printing processes."}
{"text": "U.S. Pat. No. 5,465,443 to Rice, et al., illustrates and describes various discs (and other components) of automatic swimming pool cleaners (APCs). The discs of the Rice patent typically are flexible, unitary structures defining central apertures for receiving footpads of APCs. Disc flexibility is advantageous for many reasons; flexible discs may bend when vertical or angled walls are encountered, for example, and adhere better to bottom surfaces of pools when surrounding areas are evacuated. Flexible discs additionally may more easily ride over objects extending upward from the bottom surfaces of pools as the APCs traverse those surfaces.\nMany existing discs have upper surfaces that are mostly, if not entirely, planar. This is true for discs illustrated in the Rice patent, in which only peripheral fins extend upward from the upper surfaces. Regions of the upper surfaces adjacent the central apertures, by contrast, lack upwardly-extending protrusions. See, e.g., Rice, FIG. 6.\nU.S. Patent Application Publication No. 2007/0261183 of Moore, et al., depicts additional discs and other components of APCs. Unlike discs of the Rice patent, those of the Moore application are not necessarily unitary structures, but instead may be formed of multiple parts. As illustrated, the discs may comprise forward, mid-, and rear sections. Although all sections are, to substantial extent, flexible, the mid- and rear sections beneficially are more flexible than are the forward sections. As noted in the Moore application, “Enhanced rigidity of [the] forward section additionally inhibits its assuming the shape of a corner or other transition within a pool . . . and prevents [the] forward section from folding under itself when departing from vertical surfaces such as walls.” See Moore, p. 3, col. 2, ¶ 0045 (numerals omitted). In turn, the greater flexibility of the mid- and rear sections provides improved sealing of the disc to surfaces and may improve the ability of APCs to climb pool walls. See id., ¶ 0046. However, notwithstanding enhanced rigidity of the forward section, it nevertheless may, at times, fold under itself in use.\nLike the discs of the Rice patent, those of the Moore application have generally planar upper surfaces. This is especially true for the mid- and rear sections of the discs. It likewise is true for the forward sections of the discs, although upwardly-extending peripheral fins again are shown.\nU.S. Pat. No. 6,049,933 to McLaughlin shows yet other APCs having bodies to which flexible discs may be attached via footpads. Also illustrated as attached to the bodies are bumper assemblies, which often function as leading edges of the cleaners. The exemplary bumper assemblies of the McLaughlin patent may include bumpers comprising main frame and fins. Whereas the fins are normally composed of flexible material, the main frames are substantially rigid. See McLaughlin, col. 2, ll. 50-67.\nAs depicted especially in FIG. 2 of the McLaughlin patent, no direct connection between the disc and bumper assembly exists. This is consistent with conventional designs of APCs, in which the generally planar surfaces of discs are configured intentionally to be unfettered. Indeed, past efforts of restricting (nominally vertical) movement of the generally planar surfaces have resembled the systems of U.S. Pat. No. 5,014,382 to Kallenbach, in which weight retainers separate from the discs have been used to “stop” upward movement of the discs."}
{"text": "Automatically adjustable toe angle systems for automobiles have been in existence for many years. International class B62D17 is directed to automatic camber, caster, and toe-in adjustment systems.\nWO2010045999 shows automatic adjustment devices for the toe angle of an automobile by sensing driving-critical conditions.\nEP2163412 teaches algorithms for a toe angle adjustment ECU of an automobile.\nDE 102008027326 teaches a single actuator adjustment system for simultaneously changing the toe angle of automobile wheels.\nWO2009113642 teaches a toe angle adjustment controller for dynamically adjusting the toe angle of an automobile's rear wheels.\nEP2067689 teaches a toe angle adjustment system for dynamically adjusting the rear wheel toe angle based on a steering angle velocity.\nNone of the above systems have been provided on vehicles that have an adjustable track.\nNone of the above systems are used to solve the problem of changing the vehicle track with a reduced force."}
{"text": "Traditionally, intelligent stacking technology for distributed traffic capture devices brings scalability, speed and fault tolerance to centralized network monitoring. Network traffic captured in distributed way requires a scalable centralized layer between network infrastructure and analytical equipment. This architecture creates a single logical traffic capture layer that automatically discovers and routes traffic to the designated monitoring tool via the highest-speed link and the lowest number of hops, changing the path as needed in event of link failure or reconfiguration.\nOne drawback of existing intelligent stacking technology is that various network capture devices included in the single logical traffic capture layer typically not only have different port speeds but also limited filtering resources. However, optimum path selection techniques employed by existing centralized network capture systems take into account only variable link speeds."}
{"text": "The technology relates to an exposure device including a plurality of light-emitting elements and to a light receiving device including a plurality of light-receiving elements.\nAn exposure head including a plurality of light-emitting elements such as light-emitting diode (LED) elements is used in various image forming apparatuses. Non-limiting examples of such image forming apparatuses may include an electronic printer and a facsimile apparatus that each form an image by electrophotography scheme. For example, reference is made to Japanese Unexamined Patent Application Publication No. 2007-81081."}
{"text": "One of the most basic tools which has been in use since tools were first developed is the wrench, which is a vital component of any collection of tools for mechanics and the like. A type of wrench which any mechanic will always have available is the ratchet wrench, by means of which a mechanic may rapidly and quickly accomplish either the attachment or removal of a nut on a threaded member or, install or remove a threaded bolt having a polygonal head. As is well-known, through a series of reciprocating motions of the wrench handle a ratchet wrench permits a nut to be quickly applied or removed. In the years of development of such ratchet wrenches the focus has been to increase the simplicity of such tools, the use of lightweight materials, durability and ease of use. A lightweight, low-cost, durable tool has always been the objective so as to reduce the time and effort required to perform the task for which a ratchet wrench is suitable."}
{"text": "1. Field\nThe presently disclosed subject matter relates to a vehicle reflector and a reflex pin, and particularly relates to a vehicle reflector capable of obtaining a luminous intensity in an observation angle direction which is equivalent to or higher than certain conventional reflectors by setting an inclination angle of each of the first reflection surface to the third reflection surface, and by using a reflex pin in formation of the vehicle reflector.\n2. Description of the Related Art\nIn the field of vehicle reflectors, an increase in the luminous intensity in the observation angle direction has been demanded. And, as one example of a reflector that may respond to such demand, a vehicle reflector is described in, for example, Japanese Patent No. 3340640.\nVehicle reflector 210 described in Japanese Patent No. 3340640 includes a plurality of retroreflection elements 211 for retroreflecting an incident ray Ray0 which is incident parallel with a reference axis AX in an observation angle θ direction as illustrated in FIG. 13 of the present application.\nEach of the plurality of retroreflection elements 211 includes a first reflection surface 1, a second reflection surface 2 and a third reflection surface 3 which are adjacently disposed to form a corner portion of a substantially regular hexahedron (see FIGS. 14 and 15). Each of the reflection surfaces 1 to 3 is inclined by the same inclination angle with respect to the reference axis AX (see FIG. 14), and the second reflection surface 2 and the third reflection surface 3 are disposed at the positions which are rotated by the same rotation angles (120°+β=120°6′, 120°+γ=120°6′) with the reference axis AX as the center (see FIG. 15).\nIn the vehicle reflector 210 of the above described configuration, the ray Ray0 which is irradiated from a light source 230 and parallel with the reference axis AX is incident on a plurality of retroreflection elements 211 (the first to the third reflection surfaces 1 to 3), and is reflected in the direction of the observation angle θ by the action of the plurality of retroreflection elements 211 (the first to the third reflection surfaces 1 to 3). The reflection lights Ray1 to Ray3 which are reflected in the direction of the observation angle θ respectively gather at a point P1 at the uppermost portion and points P2 and P3 at both sides under the point P1 (see FIG. 13B, and FIG. 16), and therefore, the luminous intensity in the observation angle θ direction can be enhanced."}
{"text": "1. Field of the Invention\nThe present invention relates to a device for encasing storage media wherein the protective case contains separately addressable data.\n2. Description of Related Art\nProtective cases for storage media are commonplace. Numerous diverse types of storage media, such as Compact Disc (xe2x80x9cCDxe2x80x9d), CD-Read Only Memory (xe2x80x9cCD-ROMxe2x80x9d), CD-Read/Write (xe2x80x9cCD-Rxe2x80x9d), CD-Read/Write/Re-Write (xe2x80x9cCD-RWxe2x80x9d), Digital Video Disc (xe2x80x9cDVDxe2x80x9d), Eight-Millimeter Movie Film, Video Cassette Tape and Video Game Cartridge are employed. Storage media can be easily scratched or damaged, so therefore are often sold and stored in protective cases.\nRelated art includes U.S. Pat. No. 5,119,353 to Asakura which discloses an integrated circuit chip or memory device imbedded into the central region of a CD. In U.S. Pat. No. 4,872,151 to Smith, additional data is written onto and read from the reverse side of a CD. The data described may be identification codes or possibly samples of data stored on the obverse side. U.S. Pat. No. 3,000,640 to Strauss discloses a small sample record frangibly integral with the storage case which is intended to be separated from the storage case. The Strauss invention teaches a device which must be broken from the storage case before the data on it can be read. U.S. Pat. No. 5,548,571 to Mistretta discloses a CD case which provides identification data on a storage disc by means of holes at the edge of a case. The disc is read by the disc drive while the disc is in the case.\nNone of the related art provide for the reproduction of data stored in the protective case.\nIt would be advantageous to be able to hear a message regarding the content of a CD or play a sample of music recorded on a certain CD without having to open the protective case. Similarly, it would be advantageous for a protective case to display data about a CD, such as which track is currently being played audibly from the protective case, or who the artist is. Furthermore, it would be advantageous to be able to communicate that displayed data with a computer, so it may easily be altered or read, or with a stereo amplifier, so an audio message may be amplified. Also, it would be advantageous for a protective case without a CD in it to be able to indicate audibly and/or visually which CD is to go into it.\nThe present invention contemplates a device where information may be obtained about the content of the storage media that is contained within a protective case so that the protective case is not opened unnecessarily. In particular, the present invention relates to a method and device for encasing storage media and which contains means for reproducing and/or recording messages regarding the storage media. The present invention includes the following interrelated objects, aspects and features:\nIn one embodiment, the invention comprises a protective case suitable for encasing storage media, such protective case being adaptable for placing into and retrieving such storage media; first means cooperatively attached to the protective case for storing data, such data being separate and apart from the data stored on the storage media; and second means cooperatively attached to the protective case for reproducing the separately stored data.\nIn another embodiment, the first means and the second means are attached to the interior of the protective case.\nIn still another embodiment, the first means is attached to the interior of the protective case and the second means is attached to the exterior of the protective case.\nIn yet another embodiment, the invention further comprises third means cooperatively attached to the protective case for activating the second means.\nIn still another embodiment, the third means is attached to the interior of the protective case.\nIn still another embodiment, the third means is attached to the exterior of the protective case.\nIn yet another embodiment, the invention further comprises fourth means cooperatively attached to the protective case for activating the first means.\nIn still another embodiment, the fourth means is attached to the interior of the protective case.\nIn still another embodiment, the fourth means is attached to the exterior of the protective case.\nIn still another embodiment, the invention comprises a fifth means cooperatively attached to the protective case for indicating the operation of the third means.\nIn still another embodiment, the fifth means is attached to the interior of the protective case.\nIn still another embodiment, the fifth means is attached to the exterior of the protective case.\nIn still another embodiment, the invention comprises a sixth means cooperatively attached to the protective case for indicating the operation of the fourth means.\nIn still another embodiment, the sixth means is attached to the interior of the protective case.\nIn still another embodiment, the sixth means is attached to the exterior of the protective case.\nIn yet still another embodiment, the invention further comprises seventh means cooperatively attached to the protective case for communicating the separately stored data between the protective case and remote device(s) (not shown), distant from the protective case.\nIn still another embodiment, the seventh means is attached to the interior of the protective case.\nIn still another embodiment, the seventh means is attached to the exterior of the protective case.\nIn yet another embodiment, the invention further comprises eighth means cooperatively attached to the protective case for displaying data stored in the first means.\nIn still another embodiment, the eighth means is attached to the interior of the protective case.\nIn still another embodiment, the eighth means is attached to the exterior of the protective case.\nIn yet another embodiment, the invention further comprises ninth means cooperatively attached to the protective case for storing additional data, such additional data may be displayed by the eighth means.\nIn still another embodiment, the ninth means is attached to the interior of the protective case.\nIn still another embodiment, the ninth means is attached to the exterior of the protective case.\nIn yet another embodiment, the invention further comprises a tenth means cooperatively attached to the protective case for communicating data between the protective case and remote device(s) (not shown), distant from the protective case, separate from the seventh means.\nIn still another embodiment, the tenth means is attached to the interior of the protective case.\nIn still another embodiment, the tenth means is attached to the exterior of the protective case.\nIn yet still another embodiment, the invention comprises a data storage assembly, comprising a protective case suitable for encasing storage media, such storage media comprising of component(s) from the list of: CD, CD-ROM, CD-R, CD-RW, DVD, Eight-Millimeter Movie Film, Video Cassette Tape, Video Came Cartridge and combinations thereof, such protective case being adaptable for placing into and retrieving such storage media; first means cooperatively attached to the protective case for storing data, such data being separate and apart from the data stored on the storage media, such data may be stored in a single block, and may alternatively be stored in distinct addressable segments, wherein the first means cooperatively attached to the protective case for storing data, comprising of component(s) from the list of: Data Record playback Module (xe2x80x9cDRPMxe2x80x9d), Central Processing Unit (xe2x80x9cCPUxe2x80x9d), Random Access Memory (xe2x80x9cRAMxe2x80x9d), Read Only Memory (xe2x80x9cROMxe2x80x9d), Programmable Read Only Memory (xe2x80x9cPROMxe2x80x9d), Erasable Programmable Read Only Memory (xe2x80x9cEEPROMxe2x80x9d), Electronically Erasable Programmable Read Only Memory (xe2x80x9cEEPROMxe2x80x9d) and combinations thereof, in communication with component(s) from the list of: condenser microphone, piezoelectric device and combinations thereof; second means cooperatively attached to the protective case for reproducing the separately stored data, comprising of component(s) from the list of: DRPM, CPU, RAM, ROM, PROM, EPROM and EEPROM and combinations thereof, in communication with component(s) from the list of: speaker, piezoelectric device and combinations thereof; third means cooperatively attached to the protective case for activating the second means, comprising of component(s) from the list of: contact switch, bubble switch, capacitance switch; fourth means cooperatively attached to the protective case for activating the first means, comprising of component(s) from the list of: contact switch, bubble switch, capacitance switch; a fifth means cooperatively attached to the protective case for indicating the operation of the third means, comprising of component(s) from the list of: Light Emitting Diode (xe2x80x9cLEDxe2x80x9d), Liquid Crystal Display (xe2x80x9cLCDxe2x80x9d), incandescent bulb and combinations thereof; sixth means cooperatively attached to the protective case for indicating the operation of the fourth means, comprising of component(s) from the list of: LED display, LCD display, incandescent bulb and combinations thereof; seventh means cooperatively attached to the protective case for communicating the separately stored data between the protective case and one or more remote devices (not shown), distant from the protective case, comprising of component(s) from the list of: earphone jack, stereo jack, fiberoptic interface port, infrared transmitter and detector and combinations thereof; eighth means cooperatively attached to the protective case for displaying data stored in the first means, comprising of component(s) from the list of: LED display, LCD display, incandescent bulb and combinations thereof; ninth means cooperatively attached to the protective case for storing additional data, data for use in the eighth means may also be stored in the ninth means, comprising of component(s) from the list of: DRPM, CPU, RAM, ROM, PROM, EPROM, EEPROM and combinations thereof; tenth means cooperatively attached to the protective case for communicating data between the protective case and one or more remote devices (not shown), distant from the protective case, separate from the seventh means, comprising of component(s) from the list of: earphone jack, stereo jack, fiberoptic interface port, infrared transmitter and detector and combinations thereof.\nThese and other aspects and features of the present invention will be better understood from the following detailed description of the preferred embodiments when read in conjunction with the appended figures."}
{"text": "1. Field of the Invention\nThis invention relates to a cleaning apparatus for cleaning (removing) a toner from a member to be cleaned such as an image bearing member, and further to an image forming apparatus such as a copier, a printer or a facsimile apparatus which can use such cleaning apparatus.\n2. Related Background Art\nIn a well-known image forming apparatus which repeats the step of transferring a toner image formed on the surface of an image forming member being moved to a recording material such as paper, it is difficult for all the toner to be completely transferred to the recording material during transfer, and after the transfer, part of the toner unavoidably remains on the surface of the image bearing member.\nAccordingly, it is a requisite condition for obtaining image of good quality to sufficiently remove the residual toner before entering the next image forming step, and as a cleaning apparatus for that purpose, an apparatus in which a cleaning blade formed of an elastic material such as urethane rubber closely contacts with the surface of the image bearing member to thereby remove any residual toner which has not contributed to the transfer is simple and compact in construction and is excellent in the toner removing function and has therefore been widely used heretofore.\nNow, among cleaning apparatuses, there is one designed such that collected toner is directed to a toner carrying path in the apparatus, and the toner is carried by a toner carrying screw member provided in this toner carrying path and this toner is contained and collected as a waste toner in a collected toner bottle or the like.\nIn recent years, however, it has become the mainstream to dispose a photosensitive drum and a fixing apparatus in proximity to each other to downsize the copier and shorten FCOT (first copy time) which is the first copy discharge time, and the fixing apparatus is disposed just near the cleaning apparatus, and the cleaning apparatus becomes liable to be affected by the heat of the fixing apparatus.\nAlso, from the viewpoint of using a toner of a low melting point from the necessity of curtailing electric power concerned in fixing as an energy saving countermeasure, or making the quality of image high, the development for making the particles of the toner finer has been advanced, and the following problems have arisen in the development of such a toner.\nThat is, in a cleaning apparatus 30 shown in FIG. 9 of the accompanying drawings, a screw of a spiral shape having a shaft 35a as shown in FIG. 10 of the accompanying drawings is often used as a carrying screw 35, and in this carrying screw 35, the outer diameter thereof is 15.4 mm, whereas the diameter of the shaft thereof is about 8 mm. This carrying screw 35 is rotated in the direction of arrow S in FIG. 10, whereby a waste toner is carried in the direction of arrow T.\nThere has arisen the problem that the waste toner remaining not carried by the carrying screw 35 is packed in a space surrounded by a cleaning member 34 and a scooping sheet 36 shown in FIG. 9, and this packed waste toner is heated by the heat of a fixing apparatus 40 disposed near the cleaning apparatus 30 and is crushed against an image bearing member 20, and is mechanically damaged and fused on the surface of the image bearing member 20. Particularly, the above-noted problem is liable to arise in a toner lowered in its melting point or a toner using a resin binder of low hardness to make the particles thereof finer.\nIt is an object of the present invention to provide a cleaning apparatus and an image forming apparatus which prevent the fusion of waste toner onto an image bearing member.\nIt is another object of the present invention to provide a cleaning apparatus and an image forming apparatus which removes a toner well.\nFurther objects and features of the present invention will become more fully apparent from the following detailed description when read with reference to the accompanying drawings."}
{"text": "The advantageous use of magnetic resonance technology in providing safe, rapid images of a patient has long been known. It has also been known to employ magnetic resonance technology in producing chemical shift spectra to provide information regarding the chemical content of a material.\nIn a general sense, magnetic resonance imaging involves providing bursts of radio frequency energy on a specimen positioned within a main magnetic field in order to induce responsive emission of magnetic radiation from the hydrogen nuclei or other nuclei. The emitted signal may be detected in such a manner as to provide information as to the intensity of the response and the spatial origin of the nuclei emitting the responsive magnetic resonance signal. In general, imaging may be performed in a slice or plane, in multiple planes, or in a three-dimensional volume with information corresponding to the responsively emitted magnetic radiation being received by a computer which stores the information in the form of numbers corresponding to the intensity of the signal. The pixel value may be established in the computer by employing Fourier Transformation which converts the signal amplitude as a function of time to signal amplitude as a function of frequency. The signals may be stored in the computer and may be delivered with or without enhancement to a video screen display, such as a cathode-ray tube, for example, wherein the image created by the computer output will be presented through black and white presentations varying in intensity, or through color presentations varying in hue and intensity. See, generally, U.S. Pat. No. 4,766,381.\nOne of the beneficial end uses of the disclosed systems and methods is in connection with atherosclerotic disease which is a major cause of mortality and morbidity in the United States. Localized forms of the disease, such as the deposit of plaque on the walls of blood vessels, can restrict local blood flow and require surgical intervention in some instances. While angiography is an effective means for detecting the luminal narrowing caused by plaque, it does not provide information regarding the nature of the process leading to blood flow reduction. Unfortunately, therapeutic methods, such as intravascular intervention, may experience failure due to the lack of sufficiently precise imaging methods. An imaging system capable of providing detailed, qualitative and quantitative data regarding the status of vascular walls at the time of surgical intervention, could favorably influence the outcome by enabling the selection of the intervention method to be customized to the particular need. It would also serve to provide precise guidance for various forms of localized therapy.\nIt has been known to use angioplasty and intravascular ultrasound for imaging plaques. See, generally, Spears et al., “In Vivo Coronary Angioscopy,” Journal of the American College of Cardiology, Vol. 1, pp. 1311–14 (1983); and Waller et al., “Intravascular Ultrasound: A Histological Study of Vessel During Life,” Circulation, Vol. 85, pp. 2305–10 (1992). Intravascular ultrasound, however, provides several drawbacks, including the insensitivity to soft tissue and the inability to reliably detect thrombus and discriminate thrombus (new or organized) superimposed upon plaque from soft lipid-laden plaques. Also, the presence of artifacts related to transducer angle relative to the vessel wall, and an imaging plane limited to the aperture of the transducer in variable resolution at different depths of view are further problems with this approach.\nThe feasibility of identification of atherosclerotic lesions by employing magnetic resonance (MR) microimaging in vitro has previously been suggested. See, for example, Pearlman et al., “Nuclear Magnetic Resonance Microscopy of Atheroma in Human Coronary Arteries,” Angiology, Vol. 42, pp. 726–33 (1991); Asdente et al., “Evaluation of Atherosclerotic Lesions Using NMR Microimaging,” Atherosclerosis, Vol. 80, pp. 243–53 (1990); and Merickel et al., “Identification and 3-d Quantification of Atherosclerosis Using Magnetic Resonance Imaging,” Comput. Biol. Med., Vol. 18, pp. 89–102 (1988).\nIt has also been suggested that MRI can be used for quantification of atherosclerosis. See, generally, Merickel et al., “Noninvasive Quantitative Evaluation of Atherosclerosis Using MRI and Image Analysis,” Arteriosclerosis and Thrombosis, Vol. 13, pp. 1180–86 (1993).\nYuan et al., “Techniques for High-Resolution MR Imaging of Atherosclerotic Plaques,” J. Magnetic Resonance Imaging, Vol. 4, pp. 43–49 (1994) discloses a fast spin echo MR imaging technique to image atherosclerotic plaques on an isolated vessel that has been removed by carotid endarterectomy. As the signal-to-noise ratio (SNR) decreases with the decrease in imaging time and increase in resolution, special radio frequency (RF) receiver coils were designed. The article suggests that by the use of special MR hardware at 1.5 T using various T1 and T2-weighted pulse sequences, it is possible to discriminate foam cells, fibrous plaque organized thrombus, new thrombus, loose necrosis and calcium.\nIt has also been suggested that the fat content of atherosclerotic plaque in excised tissue samples can be determined using chemical shift imaging or chemical shift spectroscopy. See, generally, Vinitski et al., “Magnetic Resonance Chemical Shift Imaging and Spectroscopy of Atherosclerotic Plaque,” Investigative Radiology, Vol. 26, pp. 703–14 (1991); Maynor et al., “Chemical Shift Imaging of Atherosclerosis at 7.0 Tesla,” Investigative Radiology, Vol. 24, pp. 52–60 (1989); and Mohiaddin et al., “Chemical Shift Magnetic Resonance Imaging of Human Atheroma,” Br. Heart J., Vol. 62, pp. 81–89 (1989).\nThe foregoing prior art articles in the aggregate could lead one skilled in the art to conclude that MR, while having potential for fully characterizing vessel wall disease, suffers from low anatomic resolution unless used in vitro on small specimens with high resolution methods.\nIt is known that in order to obtain the desired high-resolution imaging and spectroscopy of arteriosclerotic plaques, a coil can be placed close to the target blood vessel.\nIn Kantor et al., “In vivo 31P Nuclear Magnetic Resonance Measurements in Canine Heart Using a Catheter-Coil,” Circulation Research, Vol. 55, pp. 261–66 (August 1984), there is disclosed an effort to improve the SNR in the 31P spectroscopy of a dog myocardium using an elliptical coil. This coil is rigid, rather bulky, and designed for spectroscopy of the myocardium, but is not ideal for vessels.\nDisclosures of efforts to develop catheter coils for imaging vessel walls are contained in Martin et al., “MR Imaging of Blood Vessel with an Intravascular Coil,” J. Magn. Reson. Imaging, Vol. 2, pp. 421–29 (1992); and Hurst et al., “Intravascular (Catheter) NMR Receiver Probe: Preliminary Design Analysis and Application to Canine Iliofemoral Imaging,” Magn. Reson. Med., Vol. 24, pp. 343–57 (April 1992). These disclosures employ two tiny diameter, back-to-back solenoid coils to produce a good axial profile when the coils are placed along the main magnetic field.\nMartin et al., “Intravascular MR Imaging in a Porcine Animal Model,” Magn. Reson. Med., Vol. 32, pp. 224–29 (August 1994) discloses use of the system disclosed in the above-cited Martin et al. article for high-resolution images of live animals. See, also, Abstract, McDonald et al., “Performance Comparison of Several Coil Geometries for Use in Catheters,” RSNA 79th Scientific Meeting, Radiology, Vol. 189(P), p. 319 (November 1993). A strong disadvantage of these disclosures is that multislice acquisition cannot be carried out because the longitudinal coverage of the sensitive regions is limited to a few millimeters. Furthermore, the coil itself does not have the desired flexibility while maintaining the desired efficiency of data acquisition.\nU.S. Pat. No. 5,170,789 discloses a nuclear magnetic resonance (NMR) coil probe, in the form of a loop, that is said to be insertable within a specimen, which has an opening, for purposes of nuclear magnetic resonance spectroscopy (NMRS). The disclosed two component probe, which is in the nature of an endoscope to examine the colon or cervix, has a first portion which is insertable into a body cavity and a second portion which is external to such cavity. The probe has a flexible coil body with an oval or circular shape that may deform during insertion. As a result, the coil may require tuning after insertion. If the coil were made of a very rigid material, insertion problems may occur. Also, a tuning and matching circuit, in the external portion, may limit the depth of insertion.\nU.S. Pat. No. 4,932,411 discloses a probe with a transmit/receive coil for insertion in channels which are surgically or otherwise inserted in body organs, such as the brain, liver or kidneys. The coil, which is in the form of a loop, is carried and wound on the distal end of a carrier which is used to insert the coil into the body channel.\nU.S. Pat. No. 4,672,972 discloses an NMR probe disposed at the distal end of a catheter or endoscope for obtaining NMR spectra from within a patient. The multi-turn probe has a parametric amplifier and/or a gate-array attached to it and, also, has a coil cooling system.\nU.S. Pat. No. 5,413,104 discloses an invasive MRI transducer having a balloon, at least one lumen, and a flexible coil loop for insertion in a body cavity.\nIt has been known to employ an MR-active invasive device with RF transmitter coils for selective MR angiography of blood vessels. See, generally, U.S. Pat. No. 5,447,156.\nIt has also been known to employ an intravascular catheter with a Faraday screen to prevent RF electric-field interactions with the sample, such as blood, which cause the coil to detune. See, generally, U.S. Pat. No. 5,419,325.\nMR compatibility characteristics of various catheter and guide wire systems, for use in interventional MR procedures, has been considered. See Dumoulin et al., “Real-time Position Monitoring of Invasive Devices Using Magnetic Resonance,” Magnetic Resonance in Medicine, Vol. 29, pp. 411–15 (March 1993); and Abstract, Koechli et al., “Catheters and Guide Wires for Use in an Echo-planar MR Fluoroscopy System,” RSNA 79th Scientific Meeting, Radiology, Vol. 189(P), p. 319 (November 1993).\nMcKinnon et al., “Towards Visible Guidewire Antennas for Interventional MRI,” Proc. Soc. Mag. Res., Vol. 1, p. 429 (August 1994) discloses antenna designs which are asserted to make guidewires, biopsy or sample needles and other vascular interventional devices visible by MRI. One MRI stub antenna is a length of coaxial cable with 10 cm of the braid removed from the end. One end of the coaxial cable is directly connected to the surface coil input of an MRI scanner and the other end is placed in a water filled phantom. The MR image is a bright line corresponding to spins in the immediate neighborhood of the cable. A preferred MRI stub antenna is an unterminated twisted pair cable having a diameter of 0.2 or 1 mm, and a corresponding image line width of 1 or 3 mm, respectively, which provides a finer image than the coaxial cable stub antenna. A preferred combination is a steerable guidewire containing a twisted pair cable. It is suggested that a surface coil could be used simultaneously with a guidewire antenna by combining, as with phased array coils, the specimen image from the surface coil with the image of the stub antenna using the data acquired from the stub antenna, to localize the in vivo device during interventional MRI.\nIt has been known to employ an invasive device having an RF coil for transmitting RF signals which are detected by external RF receive coils to track the invasive device. See, generally, U.S. Pat. No. 5,437,277.\nIt has also been known to employ external RF transmitter/receiver coils. See, generally, U.S. Pat. No. 5,447,156.\nU.S. Pat. No. 5,323,778 discloses a probe for insertion in an artery or other body passageway. The probe has an MRI coil, an external MRI RF source and an RF heating apparatus for hyperthermia therapy.\nJapanese Kokai Patent Application No. Hei 6[1994]-70902 to Koshiichi (hereinafter “Koshiichi”) discloses two forms of dipole antenna. The first form has four wires and three ends. Two input leads are provided and two conducting poles form the dipole. One pole is inserted into a body cavity while the other is outside the body. The length of the inserted pole is about 1.2 meters at the field strength of 1.5 T common to many whole-body MRI systems today, resulting in a total antenna dipole length of about 2.4 meters. This three-ended dipole is impractical and/or has major practical disadvantages because i) the pole length is too long for applications to body cavities, blood vessels, etc.; ii) the inserted pole is loaded by the body impedance to an extent which varies with the length inserted, resulting in the MRI resonant frequency and match impedance varying with the length of insertion; iii) the fact that one end is inserted and the other end is not results in imbalancing of the impedance of the dipole which could deleteriously affect safety and performance; iv) because the length of the poles is comparable to the input leads, the tuning of the MRI resonant frequency and the impedance matching of the coil will depend on the location, orientation, and bending of each pole relative to the other; v) because the pole length is so long (about 1.2 meters at 1.5 T) this MRI probe will perform poorly compared to preexisting MRI coils whose dimensions are smaller, rendering it a disadvantageous probe design relative to pre-existing MRI probe technologies, including ones designed for use external to the body (such as solenoid “bird cage” MRI coils, surface coils, etc.); and vi) it is impractical to maintain an orientation of the dipole relative to the magnetic field when such coil is introduced into the convoluted and contorted passages of the body. For these reasons, the first dipole antenna design of Koshiichi has significant impediments for human applications.\nKoshiichi further describes a second embodiment of his probe wherein one of the poles is folded back and incorporated into a sleeve. While this may partially overcome the problem of having a free end, the other problems discussed above with respect to his first embodiment remain. In particular, the 2.4 meter length is excessively long, has inferior performance with respect to existing MRI probes which do not allow insertion, and the end with the single pole will interact with the sleeved end containing the other pole and the input leads, making it impractical to tune.\nU.S. Pat. No. 5,358,515 discloses a microwave hyperthermia applicator for limited heating of cancerous tissue including upper and lower dipole halves of the same diameter. The upper dipole half is a widened metal extension of the inner conductor of an insulated coaxial cable. The lower dipole half is a metal cylinder connected to the outer sheath of the coaxial cable. A π/2(λ/4) transformer, such as the outermost metal cylindrical sheath of a triaxial cable, is separated at its upper end from the lower dipole half which is connected to the coaxial cable outer sheath. The transformer is filled with a dielectric medium and is connected at its lower end to such coaxial cable outer sheath. When the antenna is inserted in a dissipative medium and supplied with microwave energy through the coaxial cable, only that area of the medium immediately around the antenna is heated.\nMRI has many desirable properties for the diagnosis and therapy of atherosclerotic disease. For example, it is possible to see lesions directly, even before the plaques calcify. However, the SNR of MR images obtained from conventional surface or body coils is insufficient. This is because the coils placed outside the body pick up noise from a very large region of the body. To achieve satisfactory quality, the signal receiver can be placed as close as possible to the tissue of interest (e.g., blood vessels). A coil placed on the tip of a catheter and inserted into the blood vessels could be a solution; but, the real part of the impedance of a catheter coil is relatively small and, hence, a tuning and matching circuit is preferably located immediately after the coil within the blood vessels. It is believed that prior art designs that do otherwise suffer from a significant SNR loss. On the other hand, it is believed that prior art designs, which have a tuning and matching circuit immediately after the coil in blood vessels, are too thick to be placed into small vessels.\nThere remains, therefore, a very real and substantial need for an improved apparatus and method for MR imaging and spectroscopic analysis of specimens in a manner which provides efficient data acquisition with maximum SNR while permitting in vivo or in vitro acquisition from small vessels and a wide range of other types of specimens."}
{"text": "In recent years, many varied customer needs require rapid handling of printed products consisting of signatures which are gathered for binding, trimmed, bundled for minimal shipping costs, and shipped. In a binding line, a typical operation utilizes a multiple of inserter pockets, each of which receives signatures serially from a signature supply means, opens each signature, and drops the signatures to successively straddle a gathering chain which runs in front of the inserter pockets and carries the complete collection of gathered signatures to a location for further handling to complete the binding process. Moreover, because of the need for highly efficient plant operations, there has been a constant effort to increase the speed at which machines operate which has required the development of new techniques for handling the signatures at all stages of the binding process.\nIn addition to high speed operation, it will be appreciated that any apparatus necessarily has to be compatible with the limit on the space that is available in a binding line facility. In development of the present invention, it was established as a goal for the signature feeding apparatus to address the concerns in terms of ergonomic problems, such as carpal tunnel syndrome and the like.\nCurrently, a bindery operator retrieves a small stack of signatures, such as three to four inch pile, and carries the pile to the table surface at the pocket feeder. The signature pile is compressed and gripped between the fingers and thumb, and then it is turned 90.degree. for placement on the signature backbone. The operator aligns the signatures in the small pile, fans them, and jogs them, creating a uniform series of signatures. The signatures are again compressed between the fingers and thumb and are then placed with their backbones down into the bindery feed pocket. Some strains to the operator that may occur in such loading activities include possible wrist strain when rotating the signatures, possible arm strain when lifting the pile of signatures to the table height, possible finger and thumb strain when gripping and compressing the signatures, possible wrist and arm strain when fanning the signatures, and walking fatigue in moving between a pallet having the signatures and the bindery pocket machine.\nThe assignee of this invention has been addressing the problems caused by repetitive motion in current bindery and printing tasks performed manually by operators, as can be seen from a review of its U.S. Pat. No. 5,114,129 of Chang, et al. and also U.S. Pat. No. 5,451,040 of Crabtree, and U.S. Pat. No. 5,791,643 of Bumgardner, et al. It was established in the previous patents as a goal for the signature feeding apparatus to primarily address concerns in terms of ergonomic problems, such as carpal tunnel syndrome and the like; but it was also found in solving this problem, that it was possible to increase the capacity receiving stacked signatures for feeding to the binding line within the same or a similar amount of floor space, while operating at high speed and accepting signatures in a variety of ways. The present invention is aimed mainly at the ergonomic aspects, increasing production efficiency. Also, it is preferable that the device be portable to be moved between different bindery hoppers. Alignment, jogging and aerating by bindery personnel all need to be eliminated in order to prevent the strains caused by the repetitive motions that these tasks require. The present invention addresses the needs of being able to feed signatures at high speeds to graphic arts equipment including folders, trimmers, binding lines, etc., while maintaining an ergonomically safe process and providing portability.\nThe present invention is described herein in connection with feeding financial stock or signatures which have frictional surfaces that renders them difficult to shingle into a consistent stream in a contrast to the usual catalogue or magazine stock that shingles easily into a consistent stream. Hence, the financial stock signatures described herein are printed as closed head signatures that are fed as an entire \"log\" of signatures. The present invention will be described hereinafter with respect to the formation of logs of financial signatures and feeding them into a pocket of a bindery line pocket feeder; but the present invention is not limited to this described and illustrated embodiment of the invention."}
{"text": "The present technology relates to an imaging apparatus, a focus control method, and a program, and is to be able to increase focusing precision.\nIn the related art, it is generally possible to perform an autofocus operation in an imaging apparatus. In the autofocus operation, for example, a range-finding frame which denotes which portion of a subject is to be the subject of focusing is displayed in a viewfinder, and a focus control is performed using a contrast detection method, or a phase difference detection method so that focusing is performed in the range-finding frame.\nIn addition, in Japanese Unexamined Patent Application Publication No. 2006-58405, a technology is disclosed in which focusing to a specified portion of a subject can be easily performed when performing close-up photography by switching a size of a range-finding frame between a normal mode and a macro mode. Further, in Japanese Unexamined Patent Application Publication No. 2010-237400, a technology is disclosed in which a focus detection is performed using a phase difference detection method by providing a focus detection pixel, for example, a pair of phase difference AF pixels in an imaging element."}
{"text": "Pain is a protective mechanism that allows healthy animals to avoid tissue damage and to prevent further damage to injured tissue. Nonetheless there are many conditions where pain persists beyond its usefulness, or where patients would benefit from inhibition of pain. Voltage-gated sodium channels are believed to play a critical role in pain signaling. This belief is based on the known roles of these channels in normal physiology, pathological states arising from mutations in sodium channel genes, preclinical work in animal models of disease, and the clinical usefulness of known sodium channel modulating agents (Cummins, T. R., Sheets, P. L., and Waxman, S. G., The roles of sodium channels in nociception: Implications for mechanisms of pain. Pain 131 (3), 243 (2007); England, S., Voltage-gated sodium channels: the search for subtype-selective analgesics. Expert Opin Investig Drugs 17 (12), 1849 (2008); Krafte, D. S, and Bannon, A. W., Sodium channels and nociception: recent concepts and therapeutic opportunities. Curr Opin Pharmacol 8 (1), 50 (2008)).\nVoltage-gated sodium channels (NaV's) are key biological mediators of electrical signaling. NaV's are the primary mediators of the rapid upstroke of the action potential of many excitable cell types (e.g. neurons, skeletal myocytes, cardiac myocytes), and thus are critical for the initiation of signaling in those cells (Hille, Bertil, Ion Channels of Excitable Membranes, Third ed. (Sinauer Associates, Inc., Sunderland, Mass., 2001)). Because of the role NaV's play in the initiation and propagation of neuronal signals, antagonists that reduce NaV currents can prevent or reduce neural signaling. Thus NaV channels are considered likely targets in pathologic states where reduced excitability is predicted to alleviate the clinical symptoms, such as pain, epilepsy, and some cardiac arrhythmias (Chahine, M., Chatelier, A., Babich, O., and Krupp, J. J., Voltage-gated sodium channels in neurological disorders. CNS Neurol Disord Drug Targets 7 (2), 144 (2008)).\nThe NaV's form a subfamily of the voltage-gated ion channel super-family and comprises 9 isoforms, designated NaV 1.1-NaV 1.9. The tissue localizations of the nine isoforms vary greatly. NaV 1.4 is the primary sodium channel of skeletal muscle, and NaV 1.5 is primary sodium channel of cardiac myocytes. NaV's 1.7, 1.8 and 1.9 are primarily localized to the peripheral nervous system, while NaV's 1.1, 1.2, 1.3, and 1.6 are neuronal channels found in both the central and peripheral nervous systems. The functional behaviors of the nine isoforms are similar but distinct in the specifics of their voltage-dependent and kinetic behavior (Catterall, W. A., Goldin, A. L., and Waxman, S. G., International Union of Pharmacology. XLVII. Nomenclature and structure-function relationships of voltage-gated sodium channels. Pharmacol Rev 57 (4), 397 (2005)).\nNaV channels have been identified as the primary target for some clinically useful pharmaceutical agents that reduce pain (Cummins, T. R., Sheets, P. L., and Waxman, S. G., The roles of sodium channels in nociception: Implications for mechanisms of pain. Pain 131 (3), 243 (2007)). The local anesthetic drugs such as lidocaine block pain by inhibiting NaV channels. These compounds provide excellent local pain reduction but suffer the drawback of abolishing normal acute pain and sensory inputs. Systemic administration of these compounds results in dose limiting side effects that are generally ascribed to block of neural channels in the CNS (nausea, sedation, confusion, ataxia). Cardiac side effects can also occur, and indeed these compounds are also used as class 1 anti-arrhythmics, presumably due to block of NaV1.5 channels in the heart. Other compounds that have proven effective at reducing pain have also been suggested to act by sodium channel blockade including carbamazepine, lamotragine, and tricyclic antidepressants (Soderpalm, B., Anticonvulsants: aspects of their mechanisms of action. Eur J Pain 6 Suppl A, 3 (2002); Wang, G. K., Mitchell, J., and Wang, S. Y., Block of persistent late Na+ currents by antidepressant sertraline and paroxetine. J Membr Biol 222 (2), 79 (2008)). These compounds are likewise dose limited by adverse effects similar to those seen with the local anesthetics. Antagonists that specifically block only the isoform(s) critical for nocioception are expected to have increased efficacy since the reduction of adverse effects caused by block of off-target channels should enable higher dosing and thus more complete block of target channels isoforms.\nFour NaV isoforms, NaV 1.3, 1.7, 1.8, and 1.9, have been specifically indicated as likely pain targets. NaV 1.3 is normally found in the pain sensing neurons of the dorsal root ganglia (DRG) only early in development and is lost soon after birth both in humans and in rodents. Nonetheless, nerve damaging injuries have been found to result in a return of the NaV 1.3 channels to DRG neurons and this may contribute to the abnormal pain signaling in various chronic pain conditions resulting from nerve damage (neuropathic pain). These data have led to the suggestion that pharmaceutical block of NaV 1.3 could be an effective treatment for neuropathic pain. In opposition to this idea, global genetic knockout of NaV 1.3 in mice does not prevent the development of allodynia in mouse models of neuropathic pain (Nassar, M. A. et al., Nerve injury induces robust allodynia and ectopic discharges in NaV 1.3 null mutant mice. Mol Pain 2, 33 (2006)). It remains unknown whether compensatory changes in other channels allow for normal neuropathic pain in NaV 1.3 knockout mice, though it has been reported that knockout of NaV 1.1 results in drastic upregulation of NaV 1.3. The converse effect in NaV 1.3 knockouts might explain these results.\nNaV 1.7, 1.8, and 1.9 are highly expressed in DRG neurons, including the neurons whose axons make up the C-fibers and Aδ nerve fibers that are believed to carry most pain signals from the nocioceptive terminals to the central nervous. Like NaV 1.3, NaV 1.7 expression increases after nerve injury and may contribute to neuropathic pain states. The localization of NaV 1.7, 1.8, and 1.9 in nocioceptors led to the hypothesis that reducing the sodium currents through these channels might alleviate pain. Indeed, specific interventions that reduce the levels of these channels have proven effective in animal models of pain.\nSpecific reduction of NaV 1.7 in rodents by multiple different techniques has resulted in the reduction of observable pain behaviors in model animals. Injection of a viral antisense NaV 1.7 cDNA construct greatly reduces normal pain responses due to inflammation or mechanical injury (Yeomans, D. C. et al., Decrease in inflammatory hyperalgesia by herpes vector-mediated knockdown of NaV 1.7 sodium channels in primary afferents. Hum Gene Ther 16 (2), 271 (2005)). Likewise, a genetic knockout of NaV 1.7 in a subset of nociceptor neurons reduced acute and inflammatory pain in mouse models (Nassar, M. A. et al., Nociceptor-specific gene deletion reveals a major role for NaV 1.7 (PN1) in acute and inflammatory pain. Proc Natl Acad Sci USA 101 (34), 12706 (2004)). Global knockouts of NaV 1.7 in mice lead to animals that die on the first day after birth. These mice fail to feed and this is the presumed cause of death.\nTreatments that specifically reduce NaV 1.8 channels in rodent models effectively reduce pain sensitivity. Knockdown of NaV 1.8 in rats by intrathecal injection of antisense oligodeoxynucleotides reduces neuropathic pain behaviors, while leaving acute pain sensation intact (Lai, J. et al., Inhibition of neuropathic pain by decreased expression of the tetrodotoxin-resistant sodium channel, NaV1.8. Pain 95 (1-2), 143 (2002); Porreca, F. et al., A comparison of the potential role of the tetrodotoxin-insensitive sodium channels, PN3/SNS and NaN/SNS2, in rat models of chronic pain. Proc Natl Acad Sci USA 96 (14), 7640 (1999)). Global genetic knockout of NaV 1.8 in mice or specific destruction of NaV 1.8 expressing neurons greatly reduces perception of acute mechanical, inflammatory, and visceral pain (Akopian, A. N. et al., The tetrodotoxin-resistant sodium channel SNS has a specialized function in pain pathways. Nat Neurosci 2 (6), 541 (1999); Abrahamsen, B. et al., The cell and molecular basis of mechanical, cold, and inflammatory pain. Science 321 (5889), 702 (2008); Laird, J. M., Souslova, V., Wood, J. N., and Cervero, F., Deficits in visceral pain and referred hyperalgesia in NaV 1.8 (SNS/PN3)-null mice. J Neurosci 22 (19), 8352 (2002)). In contrast to the antisense experiments in rats, genetic knockout mice appear to develop neuropathic pain behaviors normally after nerve injury (Lai, J. et al., Inhibition of neuropathic pain by decreased expression of the tetrodotoxin-resistant sodium channel, NaV1.8. Pain 95 (1-2), 143 (2002); Akopian, A. N. et al., The tetrodotoxin-resistant sodium channel SNS has a specialized function in pain pathways. Nat Neurosci 2 (6), 541 (1999); Abrahamsen, B. et al., The cell and molecular basis of mechanical, cold, and inflammatory pain. Science 321 (5889), 702 (2008); Laird, J. M., Souslova, V., Wood, J. N., and Cervero, F., Deficits in visceral pain and referred hyperalgesia in NaV 1.8 (SNS/PN3)-null mice. J Neurosci 22 (19), 8352 (2002)).\nNaV 1.9 global knock out mice have decreased sensitivity to inflammation induced pain, despite normal acute, and neuropathic pain behaviors (Amaya, F. et al., The voltage-gated sodium channel Na(v)1.9 is an effector of peripheral inflammatory pain hypersensitivity. J Neurosci 26 (50), 12852 (2006); Priest, B. T. et al., Contribution of the tetrodotoxin-resistant voltage-gated sodium channel NaV1.9 to sensory transmission and nociceptive behavior. Proc Natl Acad Sci USA 102 (26), 9382 (2005)). Spinal knockdown of NaV 1.9 had no apparent effect on pain behavior in rats (Porreca, F. et al., A comparison of the potential role of the tetrodotoxin-insensitive sodium channels, PN3/SNS and NaN/SNS2, in rat models of chronic pain. Proc Natl Acad Sci USA 96 (14), 7640 (1999)).\nThe understanding of the role of NaV channels in human physiology and pathology has been greatly advanced by the discovery and analysis of naturally occurring human mutations. NaV 1.1 and NaV 1.2 mutations result in various forms of epilepsy (Fujiwara, T., Clinical spectrum of mutations in SCN1A gene: severe myoclonic epilepsy in infancy and related epilepsies. Epilepsy Res 70 Suppl 1, S223 (2006); George, A. L., Jr., Inherited disorders of voltage-gated sodium channels. J Clin Invest 115 (8), 1990 (2005); Misra, S, N., Kahlig, K. M., and George, A. L., Jr., Impaired NaV1.2 function and reduced cell surface expression in benign familial neonatal-infantile seizures. Epilepsia 49 (9), 1535 (2008)). Mutations of the NaV 1.4 cause muscular disorders like paramyotonia congenital (Vicart, S., Sternberg, D., Fontaine, B., and Meola, G., Human skeletal muscle sodium channelopathies. Neurol Sci 26 (4), 194 (2005)). NaV 1.5 mutations result in cardiac abnormalities like Brugada Syndrome and long QT syndrome (Bennett, P. B., Yazawa, K., Makita, N., and George, A. L., Jr., Molecular mechanism for an inherited cardiac arrhythmia. Nature 376 (6542), 683 (1995); Darbar, D. et al., Cardiac sodium channel (SCN5A) variants associated with atrial fibrillation. Circulation 117 (15), 1927 (2008); Wang, Q. et al., SCN5A mutations associated with an inherited cardiac arrhythmia, long QT syndrome. Cell 80 (5), 805 (1995)).\nRecent discoveries have demonstrated that mutations in the gene that encodes the NaV 1.7 channel (SCN9A) can cause both enhanced and reduced pain syndromes. Work by Waxman's group and others have identified at least 15 mutations that result in enhanced current through NaV 1.7 and are linked to dominant congenital pain syndromes. Mutations that lower the threshold for NaV 1.7 activation cause inherited erythromelalgia (IEM). IEM patients exhibit abnormal burning pain in their extremities. Mutations that interfere with the normal inactivation properties of NaV 1.7 lead to prolonged sodium currents and cause paroxysmal extreme pain disorder (PEPD). PEPD patients exhibit periocular, perimandibular, and rectal pain symptoms that progresses throughout life (Drenth, J. P. et al., SCN9A mutations define primary erythermalgia as a neuropathic disorder of voltage gated sodium channels. J Invest Dermatol 124 (6), 1333 (2005); Estacion, M. et al., NaV 1.7 gain-of-function mutations as a continuum: A1632E displays physiological changes associated with erythromelalgia and paroxysmal extreme pain disorder mutations and produces symptoms of both disorders. J Neurosci 28 (43), 11079 (2008)).\nNaV 1.7 null mutations in human patients were recently described by several groups (Ahmad, S. et al., A stop codon mutation in SCN9A causes lack of pain sensation. Hum Mol Genet. 16 (17), 2114 (2007); Cox, J. J. et al., An SCN9A channelopathy causes congenital inability to experience pain. Nature 444 (7121), 894 (2006); Goldberg, Y. P. et al., Loss-of-function mutations in the NaV 1.7 gene underlie congenital indifference to pain in multiple human populations. Clin Genet. 71 (4), 311 (2007)). In all cases patients exhibit congenital indifference to pain. These patients report no pain under any circumstances. Many of these patients suffer dire injuries early in childhood since they do not have the protective, normal pain that helps to prevent tissue damage and develop appropriate protective behaviors. Aside from the striking loss of pain sensation and reduced or absent of smell (Goldberg, Y. P. et al., Loss-of-function mutations in the NaV 1.7 gene underlie congenital indifference to pain in multiple human populations. Clin Genet. 71 (4), 311 (2007)), these patients appear completely normal. Despite the normally high expression of NaV 1.7 in sympathetic neurons (Toledo-Aral, J. J. et al., Identification of PN1, a predominant voltage-dependent sodium channel expressed principally in peripheral neurons. Proc Natl Acad Sci USA 94 (4), 1527 (1997)) and adrenal chromafin cells (Klugbauer, N., Lacinova, L., Flockerzi, V., and Hofmann, F., Structure and functional expression of a new member of the tetrodotoxin-sensitive voltage-activated sodium channel family from human neuroendocrine cells. EMBO J. 14 (6), 1084 (1995)), these NaV 1.7-null patients show no sign of neuroendocrine or sympathetic nervous dysfunction.\nThe gain of NaV 1.7 function mutations that cause pain, coupled with the loss of NaV 1.7 function mutations that abolish pain, provide strong evidence that NaV 1.7 plays an important role in human pain signaling. The relative good health of NaV 1.7-null patients indicates that ablation of NaV 1.7 is well tolerated in these patients.\nUnfortunately, the efficacy of currently used sodium channel blockers for the disease states described above has been to a large extent limited by a number of side effects. These side effects include various CNS disturbances such as blurred vision, dizziness, nausea, and sedation as well more potentially life threatening cardiac arrhythmias and cardiac failure. Accordingly, there remains a need to develop additional Na channel antagonists, preferably those with higher potency and fewer side effects."}
{"text": "1. Field of the Invention\nThe present invention relates to a polishing method and a polishing apparatus for polishing a surface of a substrate, such as a semiconductor wafer, into a flat mirror-like surface, and also relates to a program for controlling the polishing apparatus.\n2. Description of the Related Art\nA conventionally-known polishing apparatus for polishing a substrate, such as a semiconductor wafer, comprises a polishing section, a cleaning/drying section, a transfer section, a film thickness measurement device, etc. which are unified. It is common practice in such a polishing apparatus to measure a thickness of a polishing film, present in a substrate surface, prior to polishing, and again measure the thickness of the polishing film after polishing the film and cleaning and drying the substrate. This practice is to calculate the polishing performance of the polishing apparatus by comparison of thicknesses of the polishing films before and after polishing, and to feed back the data to polishing of a later substrate.\nIn many cases, time taken for polishing a polishing film present in a substrate surface (polishing tact time), time taken for cleaning/drying the substrate after polishing (cleaning/drying tact time), and time taken for measuring the thickness of the polishing film (film thickness measurement tact time) generally differ from one another. For example, the time taken for polishing the polishing film and cleaning/drying the substrate after polishing is 60 seconds, while the time taken for measuring the thickness of the polishing film is 90 seconds. In such a case, the time taken for measuring the thickness of the polishing film is a rate-determining factor in the conventional common polishing apparatus in which a sequence of operations for a substrate is carried out successively, and the throughput of the entire polishing apparatus is limited by the film thickness measurement time.\nWhen the time taken for polishing a polishing film present in a substrate surface (polishing step) and cleaning/drying the substrate after polishing (cleaning/drying step) is shorter than the time taken for measuring the thickness of the polishing film, it is conventional practice to carry out the polishing step and the cleaning/drying step each for a predetermined tact time, and cause the substrate after cleaning/drying to stand by for the measurement of the thickness of the polishing film. In such an operating manner, a substrate, which is subject to the feedback control, must be one that is later than the substrate after polishing and cleaning/drying, waiting for the thickness measurement of the polishing film. That is, the substrate, which has undergone the predetermined tact time of polishing and cleaning/drying and is waiting for vacancy in a film thickness measurement device (and which has therefore been taken out of a cassette), is not subject to the feedback control. Therefore, if the substrate after drying, standing by for the thickness measurement of the polishing film, has been poorly polished, e.g., due to a change in the processing environment, the information cannot be reflected until the thickness of the polishing film of the substrate is measured by the film thickness measurement device. Accordingly, a substrate, which has entered into the polishing step during that period, may be poorly polished due to the change in the processing environment. In order to perform the feedback control early so as to early reflect a change in the polishing step in later substrates, it is necessary to shorten the film measurement tact time as much as possible.\nIt is also conventional practice in polishing of substrates to carry out a sequence of process steps: pre-polishing film thickness measurement step→polishing step→cleaning step→drying step→post-polishing film thickness measurement step, successively for each substrate. In this case, when the last post-polishing film thickness measurement step of the process is a rate-determining step, the upstream steps (pre-polishing film thickness measurement step, polishing step, cleaning step and drying step) must be waited for until completion of the post-polishing film thickness measurement step. Therefore, there will be a substrate which, after completion of the drying step, is standing by on a drying device and a substrate which, after completion of the cleaning step, is standing by on a cleaning device. In this case, a substrate after the polishing step cannot enter the cleaning device. Interconnects (e.g., copper interconnects) formed in a surface of a substrate, such as a semiconductor wafer, can therefore corrode after polishing, which may increase the resistance of the interconnects and provide a defective product.\nSubstrates, such as semiconductor wafers, are generally managed with a cassette, in which a plurality of substrates are housed, as a lot. All the process steps for all the substrates of one lot are not completed even when polishing and cleaning/drying of the last substrate of the lot is completed if the measurement of a thickness of a polishing film of the last substrate is not completed. Therefore, there is downtime of a polishing section during a lot change period from the termination of the polishing step for the last substrate of one lot to the initiation of polishing for the first substrate of the next lot. For the purposes of maintaining polishing environment, etc., polishing of a dummy wafer or the like in the polishing section is generally practiced during such downtime. However, polishing a dummy wafer or the like in every lot change period leads to an increased cost."}
{"text": "In recent years, work has been proceeding on fourth generation (4G) mobile communication systems. The work includes required conditions such as a required frequency band, center frequency band, coexistence with the other systems and carriers, and application of various radio environments. Requirements as a radio access technique are 100 Mbps or more in the downlink in outdoor environments, 1 Gbps in isolated cells/stationary environments, and about 1 Gbps at the maximum even in indoor environments. One of candidates for the system to adopt as a radio access technique is an Orthogonal Frequency Division Multiple Access (hereinafter, referred to as “OFDMA”) communication system using Orthogonal Frequency Division Multiple (hereinafter, referred to as “OFDM”) techniques. Furthermore, the DWCS is proposed as a mobility network configuration (Non-patent Document 1).\nFIG. 18 is a diagram showing a configuration outline of the DWCS. The DWCS is comprised of a three-layer structure of a radio access unit (hereinafter, referred to as “RAU”), base station (hereinafter, referred to as “BS”) as an intermediate layer and core network (hereinafter, referred to as “CN”). The RAU has transmission/reception antennas and signal conversion apparatus extracted from the BS that has conventionally been the lowest layer in cellular systems, thus has a simplified configuration, and performs radio transmission/reception with a mobile station (hereinafter, referred to as “MS”). The BS is connected to a plurality of RAUs through radio optical fiber cables RoF (Radio on Fiber: hereinafter, referred to as “RoF”), converts a radio signal from the RAU into a baseband signal, while converting a baseband signal to the RAU into a radio signal, and thus performs high-speed parallel signal processing. High-speed channels are connected between BSs and between the BS and CN, and constitute a mobility network.\nAs shown in FIG. 18, each RAU (for example, RAU 1, 2, 3, . . . , 10) is installed in a different position corresponding to a geographic position, service request, etc. Each BS (for example, BS 1, 2, 3) is connected to the RAU via RoF. The concept of the cell including a BS as a center in conventional cellular systems is eliminated, and as a substitute, a virtual cell (hereinafter, referred to as “VC”) is constructed using the MS (for example, MS 1, 2) as a center. For example, the MS 1 constructs a VC with RAUs 8 and 9, while the MS 2 constructs a VC with RAUs 3, 4 and 5. The mobility network control system (including the CN) selects a set of antennas i.e. a set of RAUs corresponding to a location position of an MS and service request, and allocates the set to the MS as a VC. The VC is switched according to transmission/reception radio signal environment, moving speed, location position, etc. of the MS. The DWCS enables actualization of user data transmission speeds and system capacity that are higher than in the conventional cellular system, and is considered one of promising candidates for the fourth generation (4G) mobile communication system.\nNon-patent Document 1: Shidong Zhou, Ming Zhao, Xibin Xu, Jing Wang. “Distributed Wireless Communication System: A New Architecture for Future Public Wireless Access” IEEE Communications Magazine 2003, 41 (March 3) P 108-113"}
{"text": "Electronic device components such as processors and other integrated circuits typically require a power source from which the power to operate the integrated circuit is drawn. The power source ideally provides a constant voltage, and is capable of providing a current that is sufficient to power the various devices it supplies without resulting in a drop in supplied voltage. The task of providing a constant voltage to a load that demands varying degrees of current is typically done by an electronic circuit known as a voltage regulator. Voltage regulators are typically designed to react to changes in load condition to ensure that enough power is provided to the circuit so that the voltage of the supplied power signal remains at or very near the desired level.\nVoltage regulators typically are supplied with a voltage greater than the voltage to be supplied by the regulator itself, so that the task of the voltage regulator is to use a power supply providing a large voltage and produce a power supply having a constant lower voltage. One example is a personal computer having a processor running at 3.3 volts. A power supply and voltage regulator operate together to convert the 120 volts of alternating current provided by a wall outlet power connection to a 3.3 volt direct current signal of approximately constant voltage that is used to power the processor.\nOne method of designing a voltage regulator is to change the supplied high voltage signal into a pulse width modulated signal that feeds an energy storage filter, so that the resulting output signal's voltage is dependent on the width of the high voltage pulses created in the pulse width modulation circuit. The wider the pulses are, the more energy and power is transferred to the energy storage filter circuit, and a higher voltage output signal results. Because the load on voltage regulators varies, and because changing loads will require changing current and therefore changing power be provided by the voltage regulator, the width of pulses can also be varied to change the power provided by the voltage regulator circuit at a certain voltage level. As the load varies to draw more current and therefore draw more power, the width of the pulses must be increased to provide the energy storage filter circuit with more power and therefore maintain the desired voltage at a higher current level. Circuits such as these typically therefore monitor the current drawn and the voltage level, and adjust the width of the pulses produced by the pulse width modulation circuit accordingly.\nThe frequency of the pulses created by the pulse width modulation circuit is determined by considering a number of factors, including efficiency and size of the voltage regulator circuit. Selecting too low a frequency results in components that are relatively large in physical size and that produce significant ripple or voltage variation in the output signal, while selecting a frequency that is too high results in lower overall efficiency due to the losses from the resulting large number of changes in the electronic state of the various electronic components in the voltage regulator circuit.\nA voltage regulator circuit offering minimal ripple and high efficiency that is capable of providing a constant voltage over a wide variety of load conditions is therefore desired."}
{"text": "An example fluid ejection module typically has a line or an array of nozzles with a corresponding array of flow paths. Each flow path includes a pumping chamber controllable by an associated actuator, such as a piezoelectric actuator. The piezoelectric actuator includes a layer of piezoelectric material that changes geometry (or actuates) in response to a voltage applied across the piezoelectric layer by a pair of opposing electrodes. The pair of opposing electrodes includes a drive electrode and a reference electrode, and the voltage applied across the piezoelectric layer can be controlled by maintaining the reference electrode at a fixed reference voltage or potential while varying a driving voltage or potential on the driving electrode. The piezoelectric actuator is attached to a membrane over the pumping chamber. The actuation of the piezoelectric layer in response to the voltages applied across the piezoelectric layer causes the membrane to flex. The flexing of the membrane pressurizes fluid in the pumping chamber along the flow path and eventually ejects a fluid droplet out of the nozzle."}
{"text": "With the development of Internet technologies, diverse network services can be provided for users by using an Internet service platform. Before a server corresponding to the Internet service platform provides various network services for a user, the user needs to submit login information (for example, a username and password registered with the server), so that a user identity can be verified by using the login information submitted by the user. As such, when the identity is verified, the server can provide the network service for the user, and therefore, security of the provided network service is ensured.\nHowever, with the rapidly increasing number of Internet service platforms that can provide a network service for a user, to ensure security of login information, the login information a user can register with different Internet service platforms can be the same or different, brining inconvenience for the user to maintain the login information.\nCurrently, a login information processing method is provided. To be specific, different login information is stored on a third-party server that cooperates with a server corresponding to an Internet server platform. When the user needs to use the login information, the user can obtain the needed login information from the third-party server.\nHowever, the industry needs a login information processing method that can improve login information security."}
{"text": "The following includes information that may be useful in understanding the present disclosure. It is not an admission that any of the information provided herein is prior art nor material to the presently described or claimed inventions, nor that any publication or document that is specifically or implicitly referenced is prior art."}
{"text": "Psychiatric disorders are pathological conditions of the brain characterized by identifiable symptoms that result in abnormalities in cognition, emotion, mood, or affect. These disorders may vary in severity of symptoms, duration, and functional impairment. Psychiatric disorders afflict millions of people worldwide resulting in tremendous human suffering and economic burden due to lost productivity and dependent care.\nOver the past several decades, the use of pharmacological agents to treat psychiatric disorders has greatly increased, largely due to research advances in both neuroscience and molecular biology. In addition, chemists have become increasingly sophisticated at creating chemical compounds that are more effective therapeutic agents with fewer side effects, targeted to correct the biochemical alterations that accompany mental disorders.\nYet, despite the many advances that have occurred, many psychiatric diseases remain untreated or inadequately treated with current pharmaceutical agents. In addition, many of the current agents interact with molecular targets not involved with the psychiatric disease. This indiscriminate binding can result in side effects that can greatly influence the overall outcome of therapy. In some cases the side effects are so severe that discontinuation of therapy is required.\nDepression is an affective disorder, the pathogenesis of which cannot be explained by any single cause or theory. It is characterized by a persistently low mood or diminished interests in one's surroundings, accompanied by at least one of the following symptoms: reduced energy and motivation, difficulty concentrating, altered sleep and appetite, and at times, suicidal ideation (American Psychiatric Association: Diagnostic and Statistical Manual of Mental Disorders, ed. 4. Washington, American Psychiatric Association, 1994). Major depression is associated with high rates of morbidity and mortality, with suicide rates of 10-25% (Kaplan H I, Sadock B J (eds): Synopsis of Psychiatry. Baltimore, Williams & Wilkins, 1998, p. 866). The compounds of the invention may also be used to reduce fatigue commonly associated with depression (see, for example, “Bupropion augmentation in the treatment of chronic fatigue syndrome with coexistent major depression episode” Schonfeldt-Lecuona et al., Pharmacopsychiatry 39(4):152-4, 2006; “Dysthymia: clinical picture, extent of overlap with chronic fatigue syndrome, neuropharmacological considerations, and new therapeutic vistas” Brunello et al., J. Affect. Disord. 52(1-3):275-90, 1999; “Chronic fatigue syndrome and seasonal affective disorder: comorbidity, diagnostic overlap, and implications for treatment” Terman et al., Am. J. Med. 105(3A):115S-124S, 1998).\nDepression is believed to result from dysfunction in the noradrenergic or serotonergic systems, more specifically, from a deficiency of certain neurotransmitters (NTs) at functionally important adrenergic or serotonergic receptors.\nNeurotransmitters produce their effects as a consequence of interactions with specific receptors. Neurotransmitters, including norepinephrine (NE) and/or serotonin (5-hydroxytryptamine, or 5-HT), are synthesized in brain neurons and stored in vesicles. Upon a nerve impulse, NTs are released into the synaptic cleft, where they interact with various postsynaptic receptors. Regional deficiencies in the synaptic levels of 5-HT and/or NE are believed to be involved in the etiology of depression, wakefulness, and attention.\nNorepinephrine is involved in regulating arousal, dreaming, and moods. Norepinephrine can also contribute to the regulation of blood pressure, by constricting blood vessels and increasing heart rate.\nSerotonin (5-HT) is implicated in the etiology or treatment of various disorders. The most widely studied effects of 5-HT are those on the CNS. The functions of 5-HT are numerous and include control of appetite, sleep, memory and learning, temperature regulation, mood, behavior (including sexual and hallucinogenic behavior), cardiovascular function, smooth muscle contraction, and endocrine regulation. Peripherally, 5-HT appears to play a major role in platelet homeostasis and motility of the GI tract. The actions of 5-HT are terminated by three major mechanisms: diffusion; metabolism; and reuptake. The major mechanism by which the action of 5-HT is terminated is by reuptake through presynaptic membranes. After 5-HT acts on its various postsynaptic receptors, it is removed from the synaptic cleft back into the nerve terminal through an uptake mechanism involving a specific membrane transporter in a manner similar to that of other biogenic amines. Agents that selectively inhibit this uptake increase the concentration of 5-HT at the postsynaptic receptors and have been found to be useful in treating various psychiatric disorders, particularly depression.\nApproaches to the treatment of depression over the years have involved the use of agents that increase the levels of NE and 5-HT, either by inhibiting their metabolism (e.g., monoamine oxidase inhibitors) or reuptake (e.g., tricyclic antidepressants or selective serotonin reuptake inhibitors (SSRIs)).\nThere are more than twenty approved antidepressant drugs available in the United States. The classical tricyclic antidepressants (TCAs) currently available block primarily the uptake of NE and also, to varying degrees, the uptake of 5-HT, depending on whether they are secondary or tertiary amines. Tertiary amines such as imipramine and amitriptyline are more selective inhibitors of the uptake of 5-HT than of catecholamines, compared with secondary amines such as desipramine.\nSelective serotonin reuptake inhibitors have been investigated as potential antidepressants. Fluoxetine (PROZAC®), sertraline (ZOLOFT®), and paroxetine (PAXIL®) are three examples of SSRIs currently on the U.S. market. These agents do not appear to possess greater efficacy than the TCAs, nor do they generally possess a faster onset of action; however, they do have the advantage of causing less side-effects. Of these three SSRIs, paroxetine is the most potent inhibitor of 5-HT uptake, fluoxetine the least. Sertaline is the most selective for 5-HT versus NE uptake, fluoxetine the least selective. Fluoxetine and sertraline produce active metabolites, while paroxetine is metabolized to inactive metabolites. The SSRIs, in general, affect only the uptake of serotonin and display little or no affinity for various receptor systems including muscarinic, adrenergic, dopamine, and histamine receptors.\nIn addition to treating depression, several other potential therapeutic applications for SSRIs have been investigated. They include treatment of Alzheimer's disease, aggressive behavior, premenstrual syndrome, diabetic neuropathy, chronic pain, fibromyalgia, and alcohol abuse. For example, fluoxetine is approved for the treatment of obsessive-compulsive disorder (OCD). Of particular significance is the observation that 5-HT reduces food consumption by increasing meal-induced satiety and reducing hunger, without producing the behavioral effects of abuse liability associated with amphetamine-like drugs. Thus, there is interest in the use of SSRIs in the treatment of obesity.\nVenlafaxine (EFFEXOR®) is a dual-reuptake antidepressant that differs from the classical TCAs and the SSRIs chemically and pharmacologically in that it acts as a potent inhibitor of both 5-HT and NE uptake. Neither venlafaxine nor its major metabolite have a significant affinity for adrenergic alpha-1 receptors. Venlafaxine possesses an efficacy equivalent to that of the TCAs, and a benign side effect profile similar to those of the SSRIs.\nDopamine is hypothesized to play a major role in psychosis and certain neurodegenerative diseases, such as Parkinson's disease, where a deficiency in dopaminergic neurons is believed to be the underlying pathology. Dopamine affects brain processes that control movement, emotional response, and ability to experience pleasure and pain. Regulation of DA plays a crucial role in our mental and physical health. Certain drugs increase DA concentrations by preventing DA reuptake, leaving more DA in the synapse. An example is methylphenidate (RITALIN®)), used therapeutically to treat childhood hyperkinesias and symptoms of schizophrenia. Dopamine abnormalities are believed to underlie some of the core attentional abnormalities seen in acute schizophrenics.\nA therapeutic lag is associated with the use of these drugs. Patients must take a drug for at least three (3) weeks before achieving clinically meaningful symptom relief. Furthermore, a significant number of patients do not respond to current therapies at all. For example, it is currently estimated that up to thirty percent (30%) of clinically diagnosed cases of depression are resistant to all forms of drug therapy."}
{"text": "The present invention relates to medical devices in general, and in particular to catheter ablation systems for revascularizing occluded vein grafts.\nOne of the most commonly used techniques for treating partially or totally occluded cardiac vessels is cardiac bypass surgery. With this procedure, a surgeon obtains a vessel from another portion of the patient\"\"s body and grafts the new vessel to healthy sites in the cardiac vessels in order to direct blood flow around a blockage. One of the most common vessels used in bypass surgery is a portion of the saphenous vein, which is a large superficial vein found in the leg. Such grafts are often referred to as saphenous vein grafts or SVGs.\nOne of the problems with SVGs is that they also tend to become occluded within three to five years of being grafted onto the heart muscle. For some physiological reason which is not completely understood, the material that occludes such grafts tends to be more loosely organized and brittle than the material that occludes native cardiac arteries. As a consequence, treating occluded SVGs can be more difficult because the occluding material tends to break off and can flow downstream wherein it may cause the onset of a heart attack.\nOne method of treating vein graft lesions is set forth in U.S. Pat. No. 5,681,336 to Clement et al. and assigned to the assignee of the present invention. The \"\"336 patent, which is herein incorporated by reference, discloses a system of ablating vein graft lesions including proximal and distal balloons that isolate the treatment area. In addition, the system provides for the aspiration of ablated material and/or infusion of liquids to maintain vascular pressure. In the \"\"336 patent, the ablation burrs are designed to abrade a lesion in the vein wherein the abraded material can be aspirated through a catheter that extends into the treatment area.\nWhile it is believed that the system described in the \"\"336 patent works well, additional benefits may be obtained using ablation burrs that are optimized for particle aspiration and removal of the type of blocking material found in saphenous vein grafts.\nTo improve the treatment of occluded saphenous vein grafts, the present invention comprises a system for aiding in the aspiration of ablated material from a vessel. The system includes a guide wire which is advanced into a treatment area and an ablation mechanism that is routed over the guide wire. The ablation mechanism includes an ablation burr that is rotated by a driveshaft, and a hollow sheath that extends over the driveshaft. The position of the guide wire and ablation mechanism are controlled by an advancer that moves these elements within a patient\"\"s vasculature. Rotation of the driveshaft is controlled by a prime mover, typically an electric motor or air turbine. The guide wire, driveshaft and sheath extend through a Y connector. One port of the Y connector is connected to a vacuum source that draws ablated material into a collection jar. The other port of the Y connector is coupled to the advancer.\nAccording to one aspect of the present invention, an ablation burr is designed to propel ablated material proximally into an aspiration lumen. The ablation burr includes one or more channels on the surface of the burr that direct ablated material and fluid in the vessel to the aspiration lumen as the burr is rotated. In another embodiment of the invention, the one or more channels on the burr direct fluid and ablated material and fluid in the vessel proximally and radially outward to provide a scouring action of the interior vascular wall.\nIn accordance with another aspect of the invention, an ablation burr has a proximal and distal section with the distal section having a point of maximum diameter where the proximal and distal sections meet. The diameter of the distal section tapers down to a distal tip of the burr such that the distal section is ovoidal in shape. The distal section includes one or more channels that direct ablated material and fluid towards an aspiration lumen and/or toward the interior vascular wall. The proximal section comprises a cylindrical tube of a smaller diameter than the maximum diameter of the burr. The cylindrical tube may include one or more spiral channels that direct ablated material toward an aspiration lumen. In accordance with another aspect of the invention, the distal section has a diameter that decreases linearly from the point of maximum diameter to the distal tip such that the distal section of the burr has a conical configuration.\nIn accordance with another aspect of the present invention, an ablation burr has a diameter that tapers between the point of maximum diameter and the point where the distal section of the burr joins the proximal section. This tapered section includes a number of channels that direct fluids and ablated material towards the interior vascular wall to provide a scouring action in the vessel.\nIn accordance with another aspect of the present invention, the atherectomy burr fits within a protective shroud. The atherectomy burr has a relatively flat distal face that is covered with an abrasive material. The burr has one or more tapered blades that extend proximally from the distal face that move the ablated material and liquid proximally as the burr rotates.\nIn accordance with another aspect of the invention, the ablation burr comprises an auger-type bit that cuts occluding material from the vessel and moves it proximally to an aspiration lumen.\nIn accordance with another aspect of the present invention, the ablation burr is designed as a hub of cutting blades that fit within a canister. The blades are joined at the center of the burr and extend radially outward from a central axis to the inner wall of the canister. A central lumen extends through the point at which the blades are joined so the burr can be passed over a guide wire. Each of the blades includes a tab that fits into a corresponding slot on the canister to secure the blades in the canister.\nIn accordance with yet another aspect of the present invention, the ablation burr has a xe2x80x9cdumbbellxe2x80x9d shape having proximal and distal radially expanded portions. Each of the expanded portions moves liquid and abraded material radially outward as the burr is rotated. An aperture positioned between the distal and proximal radially expanded portions is in an area of low pressure so that the aperture acts as an aspiration port to aspirate material through a driveshaft that rotates the burr.\nIn accordance with yet another aspect of the present invention, the ablation burr has a bell shape with a large central lumen that is expanded at the distal end. A distal rim of the burr is covered with an abrasive material. The central lumen of the burr allows abraded material to be gathered and directed proximally to an aspiration sheath that is positioned near the proximal end of the burr and in fluid communication with the central lumen.\nIn accordance with yet another aspect of the present invention, the ablation burr includes a series of radial holes that extend into a center lumen of the burr. A vacuum is applied to the center lumen such that occluding material in the vessels is drawn into one or more holes and is sheared off the vessel wall by rotation of the burr. The ablated material may flow through the holes and into the center lumen of the burr where it is aspirated out a center of a driveshaft or may be drawn over the burr into another aspiration lumen.\nIn accordance with another aspect of the invention, the plurality of holes are coupled to one or more fittings that mate with corresponding lumens in a catheter sheath. The additional lumens are used to aspirate ablated material drawn into the holes and to provide vacuum pressure.\nIn accordance with yet another aspect of the present invention, an ablation device comprises an outer shell having an abrasive leading surface and a core that fits within the outer shell. Liquid and material disposed between the core and the inner surface of the shell are propelled proximally by the rotation of the core and shell.\nIn accordance with yet another aspect of the present invention, the ablation burr is maintained at a fixed distance from the distal end of a sheath by a coupler having a threaded end that mates with threads on the distal end of the sheath. The ablation burr is secured to the coupler via a post having a proximal cap with a diameter that is larger than the diameter of a hole at the distal end of the coupler.\nThe post also includes a distal shaft to which the ablation burr is secured. In accordance with yet another aspect of the invention, an ablation burr includes one or more holes that eject fluid that pumped through a sealed driveshaft radially outward to scour the internal vessel walls in which the burr is being used.\nIn accordance with yet another aspect of the present invention, a driveshaft that rotates an ablation burr includes a conically shaped section at its distal end. Fluid entering the space between the conically shaped section of the driveshaft and a surrounding sheath is pushed radially outward and proximally by the rotation of the driveshaft thereby aiding in the aspiration of ablated material and liquid from a vessel.\nIn accordance with yet another aspect of the invention, an expandable sleeve fits within the sheath surrounding the driveshaft. During an ablation procedure, the sheath is extended from the distal end of the sheath and expands to seal the proximal portion of the treatment site and to aid in the aspiration of material from the vessel.\nIn accordance with yet another aspect of the invention, a method is disclosed for isolating a treatment area by routing a catheter having an inflatable balloon through a native coronary artery to the point where a bypass vessel is attached to the native artery. The balloon is inflated to seal the bypass artery so that ablation can take place without ablated material being pumped downstream.\nFinally, the present invention is a method for treating occluded native arteries by routing a catheter having an inflatable balloon through a bypass vessel and into the native artery in order to seal a treatment area such that the original blockage in the native artery can be ablated."}
{"text": "This invention relates to a process for the production of surfactant mixtures containing alkyl and/or alkenyl oligoglycoside ether carboxylic acids by reaction of an aqueous solution of alkyl and/or alkenyl oligoglycosides with ω-halocarboxylic acids, salts or esters thereof in the presence of alkali."}
{"text": "The invention relates to an automatically-closing connector for connecting a liquid injection head to an injection outlet.\nThe connector is suitable in particular for medical use and more particularly for connecting an injection head and a catheter tube or other duct of a liquid pathway to enable liquid to be inserted into the body.\nAutomatically-closing connectors are known that comprise a case which defines a housing which is accessible via an injection inlet duct opening out into the housing through an end wall of the housing, said chamber communicating with an injection outlet, and the inlet duct being designed to enable the injection head to be inserted with lateral sealing into said duct towards the housing, and the connector including a closing piston which is movable in the housing and in the inlet duct between an upstream position in which it closes the duct and towards which it is urged by resilient return means, and a downstream position in which it no longer closes the duct and towards which it is pushed by the injection head when the injection head is inserted into the inlet passage.\nAn example of such a connector is described in European patent EP 0 544 581 and the corresponding U.S. Pat. No. 5,380,306.\nWithdrawing the injection head after an injection operation can cause the patient\"\"s blood to flow back into the liquid pathway which connects the connector to the patient, thus running the risk of said pathway subsequently being blocked by coagulated blood.\nTo avoid this back flow, it is known to provide the connector with a compression chamber that is accessible to the liquid injected by the injection head and which communicates with the outlet of the connector, said chamber being designed so that its volume varies under the effect of the displacement of the closing piston so that withdrawing the injection head causes the volume of the chamber to be reduced and excess liquid to be expelled from the chamber towards the outlet, thereby establishing pressure that opposes the back flow of blood in the pathway leading to the patient.\nPublication WO 98/17192 describes an example of a connector implementing that concept.\nIn that example, the compression chamber is constituted inside the closing piston which is complex in structure.\nAn object of the invention is to provide a connector which has a compression chamber while remaining simple in structure.\nAccording to the invention, this is achieved by a connector which comprises a case which defines a housing which is accessible via an injection inlet duct opening out into the housing through an end wall (4) of the housing, said housing communicating with an injection outlet of the connector, and the inlet duct being designed to enable the injection head to be inserted with lateral sealing into said duct towards the housing, the connector including a rigid and non-deformable closing piston which is movable in the housing and in the inlet duct between an upstream position in which it closes the duct and towards which it is urged by resilient return means, and a downstream position in which it no longer closes the duct and towards which it is pushed by the injection head when the head is inserted into the inlet passage, the connector having a compression chamber accessible to the liquid injected by the injection head and which communicates with the outlet of the connector, the volume of said chamber varying under the effect of the displacement of the closing piston so that withdrawing the injection head causes the volume of the chamber to be reduced and excess liquid in the chamber to be expelled towards the injection outlet, wherein the closing piston is rigid and non-deformable, wherein said compression chamber is a space constituted in the housing around the closing piston and defined upstream by said end wall of the housing and downstream by a peripheral sealing gasket (referred to as the xe2x80x9cdownstreamxe2x80x9d gasket) carried by the piston, wherein the outlet of the connector is situated downstream from said gasket, and wherein the closing piston has an internal passage which presents an inlet that opens out into said space and which presents an outlet which communicates with the outlet of the connector.\nThe terms xe2x80x9cupstreamxe2x80x9d and xe2x80x9cdownstreamxe2x80x9d relate here to the travel direction of the fluid going from the injection head towards the outlet of the connector.\nAdvantageously, the portion of the piston that moves in the inlet duct carries a peripheral sealing gasket (referred to as the xe2x80x9cupstreamxe2x80x9d gasket) which provides sealing for the duct around the piston.\nPreferably, the closing piston has a proximal portion on which both the upstream and downstream sealing gaskets are fitted, and a distal portion of tubular shape that is directed downstream.\nThese two portions can be made together or they can be fitted one to the other. They are advantageously molded out of synthetic resin.\nSuch a piston structure is remarkably simple.\nIn a preferred embodiment, the distal tubular portion has the internal passage of the piston passing longitudinally therethrough.\nThe outlet of the connector can be constituted, for example, by a catheter tube or other liquid duct permanently fixed to the connector, or else via an outlet coupling formed on the connector and which enables a catheter or other liquid duct to be coupled thereto.\nAdvantageously, the connector has a fixed gasket placed in the housing around the piston upstream from the outlet of the internal passage of the piston and downstream from the downstream gasket of the piston so as to define a sealed space around the piston between these two gaskets, into which space the injection liquid cannot penetrate and in which it is possible to place a spring which returns the piston upstream.\nIn a particular embodiment, the fixed gasket is a part disposed so as to cover the distal end of the piston with lateral sealing of the piston, said gasket having a resiliently-opening slot such that the gasket allows the piston to pass therethrough when the piston is pushed downstream by the injection head."}
{"text": "1. Field of the Invention\nThis invention relates to an expander/prime mover for extracting heat and mechanical energy from a pressured gas, particularly to a piston type expander having dual, independent rotatable mechanical outputs.\n2. Description of the Prior Art\nThe great majority of rotary engines designed for operation by a pressured gas have utilized linearly reciprocating pistons and cylinders. Similarly, many expanders have resorted to the use of linearly reciprocating pistons and cylinders. In most instances known to Applicant, the pistons have been connected to a crankshaft and in this manner, the reciprocating movements of all the pistons produced by expansion of a pressured gas in the cylinders is converted into a rotational power output, while the gas is concurrently cooled.\nThe utilization of a crankshaft inherently involves expensive precision manufacturing operations in order to ensure the balance of the crankshaft. Additionally, the successive power strokes of each piston are delayed by an interval determined by the rotational speed of the crankshaft, since this determines the time required for the piston to return from its bottom dead center position to its top dead center position in its respective cylinder. For this reason, the utilization of fluid pressure engines is primarily confined to low torque, high speed applications. It inherently does not have the ability to generate a substantial torque at low speeds due to the substantial delay in the successive expansions of the pressured gas charges supplied to each cylinder. Additionally, there is the well known deficiency involved in every crankshaft of the effective moment arm of the piston force varying cyclically from zero to a maximum and then back to zero as the piston proceeds through its entire power stroke.\nIn co-pending application Ser. No. 634,846, filed July 26, 1984 and assigned to the assignee of this application, there is disclosed an apparatus for producing a mechanical rotating output by expansion of a pressured gas in a plurality of cylinders arranged on a circular support in an equally spaced array about the axis of a rotatable shaft to which the cylinder support is secured. The pistons cooperating with the cylinders are secured to one rotary element of a unidirectional (over-running) clutch which is constructed so that movement of the pistons in one direction relative to the cylinders is prevented while movement of the pistons in the opposite direction is permitted. A load driving element is connected to the rotatable shaft on which the cylinders are mounted and hence an intermittent rotational power output is provided by such shaft due to the reaction forces intermittently acting on the cylinders.\nIn co-pending application Ser. No. 678,439, filed Dec. 5, 1984 and assigned to the assignee of this application, a plurality of cylinders are mounted in a peripherally spaced array about the axis of a power output shaft, but the support structure for the cylinders is rigidly secured to one of the bearing supports for the shaft. The cooperating piston elements are again connected to a rotatable element of a unidirectional (over-running) clutch mounted on the shaft and impart a periodic rotation to the power output shaft through the unidirectional (over-running) clutch, but can move back to their positions relative to the cooperating cylinders free of any connection to the rotating power output shaft.\nIn both of these prior art constructions, the open inner ends of the peripherally spaced array of cylinders communicate with a fluid pressure chamber within which a small pressure is maintained, sufficient to return the piston (or the cylinders), as the case may be, to their power stroke initiating positions wherein the pistons are disposed adjacent the outer closed end of the cylinders, ready to receive another charge of pressured gas. This necessarily involves a slight loss in mechanical efficiency of the expander due to the mechanical effort required to return the relatively light weight pistons to their power stroke initiating positions in the cylinders. If the pistons can be rapidly returned to their power stroke initiating positions relative to the cylinders while developing a useful mechanical output, obviously, the overall efficiency of the expander can be improved."}
{"text": "Chemical vapor deposition involves directing one or more gases containing chemical species onto a surface of a substrate so that the reactive species react and form a deposit on the surface. For example, compound semiconductors can be formed by epitaxial growth of a semiconductor material on a substrate. The substrate typically is a crystalline material in the form of a disc, commonly referred to as a “wafer.” Compound semiconductors such as III-V semiconductors commonly are formed by growing layers of the compound semiconductor on a wafer using a source of a Group III metal and a source of a group V element. In one process, sometimes referred to as a “chloride” process, the Group III metal is provided as a volatile halide of the metal, most commonly a chlorides such as GaCl2 whereas the Group V element is provided as a hydride of the Group V element. In another process, commonly referred to as metal organic chemical vapor deposition or “MOCVD” the chemical species include one or more metal organic compounds such as alkyls of the Group III metals gallium, indium, and aluminum, and also include a source of a Group V element such as one or more of the hydrides of one or more of the Group V elements, such as NH3, AsH3, PH3 and hydrides of antimony. In these processes, the gases are reacted with one another at the surface of a wafer, such as a wafer of sapphire, Si, GaAs, InP, InAs or GaP, to form a III-V compound of the general formula InXGaYAlZNAAsBPCSbD where X+Y+Z=approximately 1, and A+B+C+D=approximately 1, and each of X, Y, Z, A, B, and C can be between 0 and 1. In some instances, bismuth may be used in place of some or all of the other Group III metals.\nIn either process, the wafer is maintained at an elevated temperature within a reaction chamber. The reactive gases, typically in admixture with inert carrier gases, are directed into the reaction chamber. Typically, the gases are at a relatively low temperature, as for example, about 50-60° C. or below, when they are introduced into the reaction chamber. As the gases reach the hot wafer, their temperature, and hence their available energy for reaction, increases.\nOne form of apparatus which has been widely employed in chemical vapor deposition includes a disc-like wafer carrier mounted within the reaction chamber for rotation about a vertical axis. The wafers are held in the carrier so that surfaces of the wafers face upwardly within the chamber. While the carrier is rotated about the axis, the reaction gases are introduced into the chamber from a flow inlet element above the carrier. The flowing gases pass downwardly toward the carrier and wafers, desirably in a laminar plug flow. As the gases approach the rotating carrier, viscous drag impels them into rotation around the axis, so that in a boundary region near the surface of the carrier, the gases flow around the axis and outwardly toward the periphery of the carrier. As the gases flow over the outer edge of the carrier, they flow downwardly toward exhaust ports disposed below the carrier. Most commonly, this process is performed with a succession of different gas compositions and, in some cases, different wafer temperatures, to deposit plural layers of semiconductor having differing compositions as required to form a desired semiconductor device. Merely by way of example, in formation of light emitting diodes (“LEDs”) and diode lasers, a multiple quantum well (“MQW”) structure can be formed by depositing layers of III-V semiconductor with different proportions of Ga and In. Each layer may be on the order of tens of Angstroms thick, i.e., a few atomic layers.\nApparatus of this type can provide a stable and orderly flow of reactive gases over the surface of the carrier and over the surface of the wafer, so that all of the wafers on the carrier, and all regions of each wafer, are exposed to substantially uniform conditions. This, in turn promotes uniform deposition of materials on the wafers. Such uniformity is important because even minor differences in the composition and thickness of the layers of material deposited on a wafer can influence the properties of the resulting devices.\nThe wafer temperature normally is set to optimize the desired deposition reaction; it is commonly above 400° C. and most typically about 700°-1100° C. It is generally desirable to operate equipment of this type at the highest chamber pressure, lowest rotation speed and lowest gas flow rate which can provide acceptable conditions. Pressures on the order of 10 to 1000 Torr, and most commonly about 100 to about 750 Torr, are commonly used. Lower flow rates are desirable to minimize waste of the expensive, high-purity reactants and also minimize the need for waste gas treatment. Lower rotation speeds minimize effects such as centrifugal forces and vibration on the wafers. Moreover, there is normally a direct relationship between rotation speed and flow rate; under given pressure and wafer temperature conditions, the flow rate required to maintain stable, orderly flow and uniform reaction conditions increases with rotation rate.\nPrior to the present invention, however, the operating conditions which could be used were significantly constrained. It would be desirable to permit lower rotation speeds and gas flows, higher operating pressures, or both, while still preserving the stable flow pattern."}
{"text": "A plethora of applications exist for effecting electrical contact between two conductors. Examples of such applications include cable connectors, PC board connectors, socket connectors, DIP carriers, etc. In an illustrative application, an interconnect system may effect an interconnection between a number of terminals on a first printed circuit board with a number of corresponding terminals on a second printed circuit board. Such apparatus are used to provide an electrical interface between two circuit boards. In another illustrative application, an interconnect system may effect an interconnection between a lead of an integrated circuit device and a conductive pad or terminal on a printed circuit board. The circuit board may then be coupled to a tester apparatus or other control means. Such apparatus are used to evaluate the performance of integrated circuit devices.\nNumerous considerations bear upon the structure of an electrical interconnect system, including both electrical and mechanical considerations. For typical interconnection systems, special attention must be given to the electrical performance thereof including self inductance, resistance, capacitance, impedance matching characteristics, etc. Mechanical considerations including life span requirements, repairability or replacability, operating temperature requirements, etc., must also be considered. Finally, specific applications of an electrical interconnect system may yield a number of unique parameters which must also be considered. For example, in an interconnect system which provides an electrical interconnection between an integrated circuit lead and a printed circuit board terminal, various parameter must be considered including the coplanarity of the terminals, the mechanical manufacturing tolerances, and the device alignment and orientation of the device terminals relative to the interconnection system.\nA main objective of an interconnection system is to maintain a non-distorting electrical interconnection between two terminals. To accomplish this, an interconnection system must be carefully designed to control the lead inductance and resistance, the lead-to-lead capacitance, the lead-to-ground capacitance, the electrical decoupling system, and the impedance matching characteristic of signal paths. All of these characteristics contribute, to some degree, to the distorting nature of the electrical interconnection system.\nVarious methods have been developed to help minimize the parasitic effects of the interrconnect system. A common method is to provide signal condition circuits adjacent the electro-mechanical contacts of the electrical interconnection system. The signal conditioning circuits, typically discrete elements such as termination components are used to adjust and control the circuit impedance. Because the requisite signal conditioning components and electromechanical contacts are physically separated, it is difficult to attain an ideal interconnect system, thereby compromising the accuracy, precision and reproducibility of the interconnect system.\nOne prior art structure is suggested in U.S. Pat. No. 4,260,762, issued on Apr. 29, 1975 to Lockhart, Jr. Lockhart suggests a test socket for interconnecting a dual-in-line integrated circuit package and a printed circuit board. A capacitor is provided in the body of the socket wherein the socket material provides the dielectric for the capacitor. The contacts of the capacitor are in contact with the socket connectors, which are in turn in contact with the integrated circuit package. That is, Lockhart suggests a test socket wherein the capacitor is provided in the socket body, rather than on the \"load board\" as previously discussed.\nA scheme to connect a first circuit board containing a test socket to a coaxial probe card, and eventually to an IC tester is suggested in U.S. Pat. No. 4,996,487, issued on Feb. 26, 1991 to Pope. The first circuit board has an integrated circuit test socket connected thereto and traces from the integrated circuit test socket to plated through-holes and further to blind vias. The coaxial probe card then engages the blind vias to provide an electrical communication path between the IC tester and the integrated circuit test socket.\nA method for reducing noise in a telephone jack is suggested in U.S. Pat. No. 4,695,115, issued on Sep. 22, 1987 to Talend. Talend suggests a modular jack for telephones in which discrete bypass capacitors are connected to the leads of the jack to filter out noise thereon. Talend contemplates using monolithic surface mount capacitors which extend to a ground plane in the modular jack element.\nThe use of a pi-network to reduce noise in a connector is suggested in U.S. Pat. No. 4,853,659, issued on Aug. 1, 1989 to Kling. Kling suggests using a planer pi-network filter comprising a pair of shunt capacitors and an inductive member in series therebetween. Kling contemplates using the pi-network filter in combination with cable connectors or the like.\nA millimeter-wave probe for use in injecting signals with frequencies above 50 GHz is suggests in U.S. Pat. No. 4,983,910, issued on Jan. 8, 1991 to Majidi-Ahy et al. In Majidi-Ahy et al. an input impedance matching section couples the energy from a low pass filter to a pair of matched, anti-parallel, beam lead diodes. These diodes generate odd numbered harmonics which are passed through the diodes by an output impedance matching network.\nFinally, a capacitively loaded probe which can be used for non-contact acquisition of both analog and digital signals is suggested in U.S. Pat. No. 5,274,336, issued on Dec. 28, 1993 to Crook et al. In Crook et al., the probe consists of a shielded probe tip, a probe body which is mechanically coupled to the probe tip, and an amplifier circuit disposed within the probe body."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of golf. In particular, the invention is concerned with a system for changing the trajectory of the golf ball when struck by the face of a golf club by applying a coating friction altering material to the club face.\n2. Description of the Prior Art\nAs golfers are well aware, an unwanted hook or slice when hitting the golf ball can be detrimental to the enjoyment of the game. A hook or slice is generally the result of side spin imparted to the ball during impact. This spin occurs when the club face is not perpendicular to the arc of the swing. The greater the deviation from perpendicular, the greater the resulting spin and resulting hook or slice. The prior art has addressed this problem by changing the surface characteristics of the golf ball. This technique has produced mixed results."}
{"text": "Fluid pressure operated devices for radially expanding portions of tubing are known, such devices being disclosed in U.S. Pat. Nos. 1,448,457; 2,479,702; 2,938,562; and 3,200,627. Structures of these patents include fluid conveying mandrels receivable within the tubes, means for expanding the mandrels to hold the tubing, and die elements against which portions of the tubes are flared or otherwise formed by fluid pressure. These expanding devices are customarily used to provide shapped bulges in tubing for coupling purposes and the like. In addition, they are used to expand tubing to fit fins and other apparatus connected to the outer surface of the tubing.\nSeveral hydraulic chucks have been developed which are suitable for tightly gripping one end of the elongated tube while simultaneously delivering a tube expanding fluid. Among these U.S. Pat. Nos. 3,505,846; 3,813,751; 3,962,769; and 4,189,162.\nOne of the problems which continuously has plagued those users of hydraulic chucks is the inability of such chucks to accomodate relatively large diameter tubes such as those having diameters of from about 2 inches up to about 7 or 8 inches or more. While many of the designs are suitable for gripping the tube satisfactorily, it has not been possible to adequately seal the device. On the other hand, where the sealing may be effective, the device is inadequate for gripping the tube to hold it in position under the substantially high pressure delivered by the hydraulic fluid. In addition, many difficulties with present hydraulic chucks are incurred because the piston which is driven by the fluid to cause activation of the gripping means such as radial collets, forces the piston against the tube. In large diameter operations, the force required to cause the collets to adequately grip the tube often times causes damage to the tube or relative movement of the tube with respect to the collets. In any event, the damage to the tube renders the operation substantially more expensive if the tubes can be salvaged. Repair, or the use of substantially lower pressures prevent wide acceptance of hydraulic chucks for larger diameter tubes.\nAccordingly, a hydraulic chuck would be highly desireable that could be employed with large diameter tubes of at least 2 to 8 inches in diameter or larger. The device must be capable of gripping the tube firmly to prevent relative movement thereof during application of hydraulic pressure, while simultaneously be capable of gripping the tube without damaging the tube."}
{"text": "The present invention relates to methods and apparatus for data transmission. More particularly, the invention relates to architectures and techniques for implementing a common transceiver design that enables data transmission in any of the three standards: 100BaseT2, 100BaseTX, and 100BaseT4.\nIn modern data networks, e.g., computer networks, there exist multiple standards for data transmission. There are, for example, currently three standards for transmission of 100 Megabits per second (Mbps) over copper twisted pairs up to a length of 100 meters in an Ethernet Carrier Sense Multiple Access with Collision Detection (CSMA/CD) local area network. These standards are 100BaseT2, 100BaseT4, and 100BaseTX. Each of these standards has its advantages and disadvantages. For example, the 100BaseTX standard permits the use of lower cost circuitry over more expensive wiring while the 100BaseT2 standard requires more complex and thus expensive circuitry over cheaper wiring. The 100BaseT4 standard is generally thought of as being somewhere in the middle. The physical level interface for each standard is thus different, and therefore a physical device (PHY) that supports one of the standards is not inter-operable with a device that supports another standard.\nAt the start of a channel link, a common protocol, termed \"auto-negotiation,\" is initiated in order to identify the nature of the standard that is supported by the device at the other end of the transmission channel. By way of example, auto-negotiation may be performed between a network \"hub\" and the network interface card (NIC) of a terminal, which is attempting to communicate with the hub, to ascertain the standard with which the network interface circuitry communicates. If the devices at opposite ends of the channel are not compatible, a link is either not possible, or more practically, the link may be \"negotiated\" down to a 10BaseT interface, which is a 10 Mbps link.\nFIGS. 1 and 2 are prior art illustrations showing respectively the transmission architecture associated with the 100BaseT2 standard and a typical transceiver architecture therefor. Referring to FIG. 1, the 100BaseT2 standard typically employs only two pairs of copper wires, 102 and 104. Over each pair, data transmission takes place at full-duplex, with a bit rate of about 50 Mbps, a baud rate of about 25 Mega bauds per second (MBps). The analog-to-digital (A/D) sampling rate (e.g., on A/D 106 or 108) is about 75 Mega samples per second (Msamples/sec), as is the digital-to-analog (DAC) sampling rate (e.g., on DAC 110 or 112). The sampling rate, as is known, may represent any integer multiple of the baud rate.\nFIG. 2 shows the 100BaseT2 transceiver architecture from one end of the channel link. The 100BaseT2 standard employs digital processing techniques, implemented by DSP block 114 in FIG. 2, to compensate for, among others, cable distortion and to perform cross-talk suppression. For each of the full-duplex twisted pair links, the physical device (PHY) dedicates one Digital to Analog Converter (DAC), one transmit filter (Tx filter), one Analog to Digital Converter (A/D), and one receive filter (Rx filter). The transmit and receive filter pairs (Tx/Rx filter pairs) for channels 102 and 104 are shown in FIG. 2 as Tx 116/Rx 118 pair and Tx 120/Rx 122 pair, respectively. As mentioned earlier, the A/Ds and DCs typically sample at 75 MHz, and the signals in the analog filters contain components up to around 25 MHz. In the DSP block the device may also contain two Echo/Near End Cross-Talk Cancellers and Two Equalizers (conventional and not shown in FIG. 2). This digital circuitry also operates typically at 75 MHz. Transformers 124 and 126 are employed to conform the data to the requirements of the standard, which is a conventional technique. Note that a 100BaseT2 transceiver has two of each analog circuit elements required to perform receive and transmit.\nFIGS. 3 and 4 are prior art illustrations showing respectively the transmission architecture of the 100BaseTX standard and a typical transceiver architecture therefor. Referring to FIG. 3, the 100BaseTX standard also employs only two pairs of copper wires, 132 and 134. Over each pair, data transmission takes place at only half-duplex, with one pair carry data from master to slave and the other from slave to master. The bit rate is about 100 Mbps, and the baud rate is typically about 125 Mega bauds per second (MBps), which includes coding required for DC control and other signaling functions. The A/D sampling rate (e.g., on A/D 136) is typically at the baud rate, i.e., about 125 Msamples/sec, as is the DAC sampling rate (e.g., on DAC 140).\nFIG. 4 shows the 100BaseTX transceiver architecture from one end of the channel link. The 100BaseTX standard also employs digital processing techniques, illustrated in FIG. 4 as DSP block 144, to also compensate for, among others, cable distortion. Since the environment in which the 100BaseTX standard is employed is typically less harsh than that associated with the 100BaseT2 environment, e.g., there is typically lower cross talk and little, if any, echo, the set of DSP functions in DSP block 144 is typically less complex than those implemented in DSP block 114 of the 100BaseT2 architecture. However, the higher sampling rate (i.e., 125 Msamples/sec versus 75 Msamples/sec) requires higher performance A/Ds and DACs in the conventional 100BaseTX implementation.\nFIGS. 5 and 6 are prior art illustrations showing respectively the transmission architecture of the 100BaseT4 standard and a typical transceiver architecture therefor. Referring to FIG. 5, the 100BaseT4 standard employs four pairs of copper wires, 160, 162, 164, and 166. On pairs 164 and 166, the transmission is full duplex, and is half duplex on pairs 160 and 162. The bit rate is about 33.33 Mbps, and the baud rate is typically about 25 Mega bauds per second (MBps). The A/D sampling rate (e.g., on A/D 170, 172, or 174) is typically at 50 Mega samples/sec, as is the DAC sampling rate (e.g., on DAC 176, 178, or 180).\nFIG. 6 shows the 100BaseT4 transceiver architecture from one end of the channel link. The 100BaseT4 standard also employs digital processing techniques, implemented in FIG. 6 by DSP block 182, to also compensate for, among others, cable distortion and to perform cross-talk suppression. The 100BaseT4 standard employs more pairs of wires at a lower baud rate than the 100BaseT2 standard, and the impairments are less severe than those seen in the 100BaseT2 implementation. Consequently, the set of DSP functions in DSP block 182 is typically less complex than those implemented in DSP block 114 of the 100BaseT2 architecture.\nIn the prior art, a transceiver is typically custom designed to transmit and receive in only one of the three standards. In other words, a given transceiver, such as that illustrated in one of FIGS. 2, 4, and 6, is not inter-operable with devices, hubs, or network interface cards (NIC) that adhere to a different standard or multiple standards. Consequently, it is typically necessary to design, manufacture, stock, and support different transceivers for different standards. Additionally, when transceivers are implemented as integrated circuits, the prior art custom approach involves a substantial chip die size overhead. By way of example, to implement 100BaseT2 and 100BaseTX, the prior art custom approach would require 50% more die size, and in order to implement 100BaseT4 with 100BaseT2 and/or 100Base TX, the prior art would require 80% to 100% more die size. As can be appreciated, this approach involves a substantial amount of design and manufacturing effort and expense.\nIn view of the foregoing, there is desired a common denominator transceiver architecture and methods therefor that support all three 100 Mbps standards: 100BaseT2, 100BaseTX, and 100BaseT4. To reduce manufacturing and implementation costs, the common denominator transceiver architecture and methods therefor preferably allow, with minor modifications, circuits of the common denominator transceiver architecture to be reused when implemented to support different standards."}
{"text": "Neuropathic pain is a complex, chronic pain state that is generally accompanied by tissue injury. With neuropathic pain, the nerve fibers themselves might be damaged, dysfunctional, or injured. These damaged nerve fibers send incorrect signals to other pain centers. The clinical causes of neuropathic pain are widespread and include both trauma and disease. For example, traumatic nerve compression or crush and traumatic injury to the brain or spinal cord are common causes of neuropathic pain. Furthermore, most traumatic nerve injuries also cause the formation of neuromas, in which pain occurs as a result of aberrant nerve regeneration. In addition, cancer-related neuropathic pain is caused when tumor growth painfully compresses adjacent nerves, brain or spinal cord. Neuropathic pain is associated with diseases such as diabetes or alcoholism.\nUnfortunately, the available drug therapies for neuropathic pain often do not provide the patients in need of such treatment the adequate pain relief. In addition, current therapies have serious side-effects including, for example, cognitive changes, sedation, nausea and, in the case of narcotic drugs, addiction. Many patients suffering from neuropathic pain are elderly or have other medical conditions that particularly limit their tolerance of the side-effects associated with available drug therapy. A number of anti-inflammatory, anxiolytic, narcotic and even anti-convulsants are currently used by the practitioners to treat neuropathic pain, but with limited success.\nThe inadequacy of current therapy in relieving neuropathic pain calls for new compositions and methodologies of addressing the physical and social needs of the patient suffering from such condition. Methods of alleviating neuropathic pain would improve the quality of life for many people suffering from pain due to trauma or disease."}
{"text": "ANSI/IEEE Std 802.11, 1999 Edition 802.11 standard is a standard instituted by the Wireless Local Area Network Standard Working Group for physical layer protocol and medium access control (MAC) protocol, which facilitates the establishment of the interoperating network equipment by the wireless local area network devices manufacturers and wireless devices manufacturers. Physical layer defines the signal feature and signal modulation of data transfer. The medium access control (MAC) layer can be deemed as being composed of a series of services which can accomplish the functions such as information exchange, power control, association management, synchronal management and process management.\nWLAN (wireless local area network) means interconnecting computer devices by using wireless communication technology, so as to configure a network system capable of communicating with each other and implementing a resource sharing. In a WLAN network environment, an extended service set system is composed of one or a plurality of base service sets and a distribution system in the connection with the above, each of which at least comprises one wireless access point apparatus. The extended services sets are distinguished by ESSID (Extended Services Sets Identifier) and the wireless access points are distinguished by BSSID (Basic Services Sets Identifier). A station obtains the information of the Extended Services Sets and Basic Services Sets in the present WLAN environment by scanning channels and selects to join an appropriate basic services set. Load balance in WLAN is a strategy, under which a wireless local area network system comprises a plurality of access points providing access services to a variety of mobile stations, wherein each of the access point members allows or rejects the access of a mobile station by station management using network information, thereby balancing the network load, taking full advantage of the available network resource and reducing the network congestion. Usually the distribution of access points joining the load balancing should meet the conditions that a station can selectively get access to another access point with lower load, when the access point rejects the access of the station owing to overload; access points meeting the above-mentioned conditions are divided in a group called load group, the station must selectively get access to an access point in the load group, the access point having relative lower load and relative most appropriate for the access of the station.\nThe current wireless local area network protocol has not stipulated how the load balance processing is conducted and how the load balance information is exchanged in the network. The regular implementing method is collecting and managing the network load balance information under the information exchange protocol between the access points, wherein the load balance information brings a great deal of additional overhead when it is being broadcasted in the network; moreover, the unicity of the load balance information renders a low quality of the station access, which makes against the optimized management to the whole network. In regard to a WLAN network conducting access and resource optimization, it is advantageous for the improvement of the quality of the network communication to implement an optimized load balancing system with high performance."}
{"text": "In a variety of document presentation systems such as printing systems, it is common to rasterize data to generate a bitmap representation of each sheetside image of the document by processing a sequence of data objects. The data objects are typically initially defined in a page description language or other suitable encoding and at some point prior to writing to a bitmap are represented as regions of rectangles of pixels. Typically, the sheetside image is then generated into a bitmap memory as a two dimensional matrix of pixels representing the intended document sheetside image, and subsequently compressed.\nOften these sheetside images are saved to enable reuse during printing of a job. Upon reuse, the sheetside images are decompressed, further processed and re-compressed. However, the decompression and re-compression processes incur inefficiencies.\nAccordingly, an efficient processing mechanism is desired."}
{"text": "1. Field of the Invention\nThe present invention relates to highly-oriented diamond films of which crystal grains are highly oriented, methods for manufacturing the same, and electronic devices having highly-oriented diamond films. Specifically, the present invention relates to highly-oriented diamond films suitable for electronic devices such as transistors and diodes, methods for manufacturing the same, and electronic devices having the highly-oriented diamond films.\n2. Description of the Related Art\nA highly-oriented diamond film is a polycrystalline film in a broad sense, but a growth direction and an in-plane direction of the crystal grains are each oriented in one direction and the surface is in a characteristic state of a series of flat (001) facets. Therefore, the highly-oriented diamond film is known that the crystal defect density near the surface is lower than that of general polycrystalline films and that the carrier mobility near the surface is about two figures higher than that of general polycrystalline films. Because of this, it is thought that the highly-oriented diamond film is suitable for electronic devices such as a field-effect transistor utilizing the carrier mobility in the lateral direction.\nSuch a highly-oriented diamond film can be formed by, for example, subjecting a silicon substrate to microwave radiation while applying a negative bias voltage in a gas phase containing methane gas (U.S. Pat. No. 5,523,160). Furthermore, a method for preparing a heteroepitaxial diamond film having a flat surface is suggested in B. Dischler, C. Wild, eds. Low-Pressure Synthetic Diamond, 1998:153-158 (referred to as non-patent document, hereinafter). In the method, after the application of a bias voltage to a substrate, predominant [001] growth of diamond (a first growth) is performed at a rate of 3.9 to 4.5 μm/hr for improving a degree of orientation of crystalgrains, and then surface-planarizing growth (a second growth) of diamond is performed at a rate of 2 μm/hr or less for planarizing the surface.\nThe above-mentioned conventional technologies have some problems as follows: In an electronic device such as a transistor having a diamond film, properties of the device are improved with an increase in the diamond crystal grain size. For example, in the application to a transistor, since a gate length is generally 30 to 100 μm, the highly-oriented diamond film must have a crystal grain size of at least 30 μm, preferably 100 μm or more. When the highly-oriented diamond film is prepared according to the method disclosed in the non-patent document, diamond crystals grow in a columnar form in the first growth to gradually increase the grain sizes. However, since the increasing rate is small, a film having a large thickness is required for ensuring a crystal grain size of 30 μm or more; which is a problem.\nFurthermore, a diamond film having a flat surface is desirable. However, in the method described in the non-patent document, the diamond in the second growth is synthesized under such conditions that a growth rate in [111] orientation is higher than that in [001] orientation. Therefore, even if protruding non-oriented crystals are slight, the non-oriented crystals predominantly grow. When a highly-oriented diamond film containing non-oriented crystals is applied to an electronic device as a substrate, the surface is planarized by polishing after a film is epitaxially deposited on the diamond film for preventing inner damage caused by the polishing and for doping. This causes changes in growth sectors at the non-oriented crystal portions. Since distribution coefficients of an impurity differ in the growth sectors, the non-oriented crystal portions and the other portions each have an impurity content being different from that of each other when the epitaxially grown film is doped with the impurity. This is undesirable for device movement. Therefore, when the method disclosed in the non-patent document is conducted, the first growth must be performed until the non-oriented crystals do not protrude. Consequently, a thin film, for example, having a thickness of 50 μm or less cannot be formed by this method; which is a problem.\nThe non-patent document discloses that a smooth diamond film having a film thickness is 5 μm or less can be formed by using a favorable SiC film as a buffer layer. However, in this case, a ratio of the average crystal grain size to the film thickness (average crystal grain size/film thickness) is still ½ or lower, so a large film thickness is required for yielding an average crystal grain size of 30 μm or more. This problem has not been solved yet."}
{"text": "Semi-automatic pistols can be divided into a number of different types. These include pistols that use a blowback mechanism and those that utilize a short-recoil Browning-type mechanism.\nWith blowback pistols, only the slide moves relative to the frame of the gun upon discharge. The barrel of a blowback pistol is fixed to the frame but the slide is not secured to the barrel. Immediately after firing the pistol, the recoil force starts to drive the slide rearward commencing the extraction of the spent cartridge case. Accordingly, part of the extraction of the case occurs during the high-pressure period of the firing cycle. If the slide is too light the case is extracted too soon affecting the loading cycle. The blowback mechanism is typically used with low-powered cartridges. Semi-automatic pistols for higher-powered cartridges generally utilize a short-recoil mechanism.\nWith short-recoil operated pistols, both the barrel and slide move together rearward upon discharge of the gun. Prior to the firing of the cartridge, the barrel is engaged to the slide by a locking mechanism. After firing, the recoil force drives both the slide and barrel rearward, but since they are in engagement, the extraction of the case has not started. After the high-pressure period has passed, an actuator begins to disengage the barrel from the slide. The barrel travels a short distance before coming to rest forward of the magazine, hence short-recoil, and is completely disengaged from the slide. The slide continues and begins extraction of the spent cartridge case using its kinetic energy and the residual gas pressure in the barrel. After extraction, the spent case is ejected. The slide continues until full travel is reached.\nIn short-recoil operated pistols, the barrel may be locked to the slide by a number of locking mechanisms. The barrel may be provided with peripheral ribs, studs, lugs or other mechanism and may be rotated, cammed or otherwise engaged and disengaged from the slide. Alternatively, a separate locking block may be used to lock the barrel to the slide. A common method utilizing a locking block is the dropping barrel method as depicted in U.S. Pat. No. 4,915,011 hereby incorporated by reference in its entirety.\nIn the dropping barrel method of locking, the barrel is slidably mounted for straight line longitudinal motion, and the locking mechanism comprises a separate locking block provided with an actuation mechanism for engaging the barrel to the slide. The actuation mechanism comprises a cam in the frame to drive the barrel downwardly from engagement with the slide during initial rearward movement of the slide from the battery position and upwardly to engage the barrel to the slide during final movement of the slide to the forward battery position.\nIn metal frame pistols, the slide is usually secured for such movement by longitudinally spaced pairs of metal guide rails. The guide rails generally include four rails, one forward and one rearward pair.\nIn recent years there has been a trend in the firearm industry to utilize polymers in the manufacture of semi-automatic pistols, particularly in fabricating unitary frames therefor by injection molding techniques. Generally, in such frames, the front pair and rear pair or spaced guide rails are partially embedded in the polymer of the frame. The cam, which is part of the locking block, is assembled into the frame and held by cross pins. The locking block is located between the two sets of rails. Polymer frame pistols are desirable in that they are lighter than pistols with metal frames.\nCurrently, firearm manufacturers are making compact polymer frame pistols which feature a reduced length frame, slide and barrel as disclosed in U.S. Pat. No. 5,717,156, which is hereby incorporated by reference in its entirety. Compact pistols are smaller and lighter than standard size semi-automatic pistols. Additionally, the reduction in trigger reach and grip circumference of compact designs increases concealability and is thought to enhance shooting ergonomics. Compact designs are available for both low-powered and high-powered cartridges.\nIn some prior art compact pistols, there is a limit on the reduction in size of the pistols. The size limit is due to the fact that the guide rails and locking block are separate components. Given the separation of guide rails and locking block, and the configuration of prior art frames, a minimum distance between the guide rails and locking block is necessary."}
{"text": "The present disclosure relates generally to cooling of electronic components, and more particularly to cooling of such components using piezoelectrically driven structures.\nThere is a general trend in the electronics industry toward smaller and higher power integrated circuits and electronic devices. As the operating frequency and number of transistors on each die increases, so does the heat produced. For example, in the past few years the power dissipated from central processing units of laptop computers has gone up from 4 W (e.g., a 233 MHz Intel Pentium™ processor with MMX) to as much as 22 W (e.g., a 1.2 GHz Intel Pentium III-M™ processor). Newer central processing units are using and dissipating even more energy. Heat generation and heat density in electronics will continue to increase as more processing power is desired in laptop computers, web servers, cellular telephones, cellular phone base stations, PDAs, and other electronic devices. Portable electronic devices will require more cooling and devices that have up to now relied on natural convection alone for cooling may henceforth require active cooling elements.\nConventional arrangements for cooling electronics include axial fans and heat sinks, either alone or in combination. Axial fans drive air by rotating a set of fan blades that are mostly perpendicular to the axis of rotation. The operating lifetime of an axial fan is directly related to the life expectancy of its bearings. Bearings that are smaller than those typically used in 25×25×12 mm axial fans do not last long. Thus, reducing the size of a fan and its bearings dramatically decreases the life expectancy of the fan. Consequently, axial fans are generally not reliable when made smaller than 25×25×12 mm. It is therefore problematic to include axial fans in many portable electronic devices.\nPassive cooling systems that have no moving parts are often used in electronic devices because of the long lifetime of passive cooling systems. Heat sinks and heat pipes fall into this category. Heat sinks rely on natural convection over a large surface area to dissipate energy to ambient air. If the amount of energy to be dissipated is increased, the heat sink must have more surface area and therefore may need to be increased in size. This can make heat sinks undesirable when space is at a premium. Heat sinks therefore may not be suitable for cooling hot electronics in tight spaces. Heat pipes may function to move heat away from a local source, but heat pipes, like heat sinks, require large surface areas to dissipate heat through natural convection at a radiator. Moreover, heat pipes generally do not work well in small and confined spaces. Conventional cooling technologies typically do not operate well or reliably when miniaturized for use importable electronics.\nIt has previously been proposed to construct fans by attaching a piezoelectric actuator to a fan blade to drive ambient air. Most of these designs use at least one substantially elongated, flat, cantilevered blade. The blade or blades are vibrated, typically by using a bending type piezoelectric element at or near the first bending mode natural frequency of the blade and piezoelectric element system. When two blades are used, they are vibrated out of phase with each other to reduce vibrations in the housing and mounting system. Such piezoelectric fan designs use low frequencies (under 400 Hz) to obtain large amplitudes at the free end of the cantilever blade. Also most move air by vortex shedding due to large vibration amplitudes at the free end of the blade or blades. Voltage requirements (over 100 V), fan blade displacement (over one inch displacement at the fan blade end) and cost have generally been too great for the piezoelectric cooling systems to be used successfully in portable electronics. Furthermore, the piezoelectric cooling devices cannot easily be scaled down because the natural frequencies of the blades increase as the size decreases, and the bending-type piezoelectric elements tend to perform poorly at the higher frequencies that are required.\nIt would be desirable to provide a cooling system that has one or more of the following characteristics: (a) be able to dissipate a large amount of heat in a small volume; and (b) fit easily into portable and non-portable electronic systems that have high packaging density. It would be further desirable that such a system be inexpensive, efficient, robust, operable from batteries and/or easy to design into a system. For the latter characteristic, it would be desirable that the system use surface mount technology."}
{"text": "As communications devices become more capable of handling multiple tasks, there is less of a need for people to carry as many items with them on-the-go, leading to a minimization trend. This trend is particularly apparent in the transition from simple cell phones to newer generations of cell phones integrating the calling function with e-mail, mobile apps, movie viewing, GPS navigation, calendar and clock functions, etc. In addition, the sizes and shapes of modern cell phones are much more compact than in the past, now approaching the size of other items commonly carried on-the-go, such as credit cards, driver licenses, paper money, etc.\nTo facilitate convenient carrying of this smaller number of items, there is a need for a single carrying device which is capable of carrying both the mobile communications devices (e.g., cellular phones) or PDAs, along with the carrying of these miscellaneous other items, such as money and credit-card sized items in a single convenient pocket case."}
{"text": "Overhead storage bins in aircraft typically include latching systems to secure the pivotal bucket to the upper housing. However, many latch systems have proven unreliable. Accordingly, a need exists for an improved latch system. U.S. Patent Publication No. 2014/0197721 and U.S. patent application Ser. No. 14/796,829, filed Jul. 10, 2015 are incorporated by reference herein in their entireties."}
{"text": "This invention relates generally to portable drawing surfaces and more particularly to a base structure with interchangeable components. The invention permits the manufacture of a base structure that can be combined with an accessory unit and a working surface to create either a lighted drawing surface or a drawing surface with a utensil accessory tray. Without this invention, it would be necessary to manufacture separate base structures to create drawing surfaces of different characteristics. For example, one manufacturing line would be needed for molding base structures for a drawing surface with a utensil accessory tray. An entirely separate line with another set of manufacturing equipment and operating personnel would be needed to mold base structures for a lighted drawing surface. Thus, a great deal of wasteful duplication in machinery and personnel would be necessary to make two distinct drawing surfaces.\nThis invention obviates the need for such duplicative efforts. The interchangeable nature of the base structure and other components enables the manufacture of multiple drawing surfaces with distinctive characteristics from the same base parts. This invention greatly reduces molding costs since molds for only one base structure, not two, are needed. This invention results in substantial savings in manufacturing costs as tooling requirements are substantially lowered. Thus, this invention permits the economical manufacture of portable drawing surfaces with either an accessory tray for storing drawing utensils or a backlit working surface."}
{"text": "1. Field of the Invention\nThe present invention relates to a reconfigurable toy assembly, and more particularly, to a reconfigurable toy assembly wherein a toy vehicle, formed by separatably combining a tractive unit and a trailer section, is constructed such that it is possible to reversibly reconfigure the tractive unit and the trailer section into a toy robotic humanoid and a play space, respectively.\n2. Description of the Prior Art\nSuch reconfigurable toy assemblies of this type have recently been proposed and put into practical use that toy cars simulating a bus, truck and the like are reconfigured into a toy house, toy robotic humanoid and the like. However, any of these reconfigurable toy assemblies can be reconfigured into only one kind of thing; hence, the play with such a reconfigurable toy assembly may be apt to be monotonous."}
{"text": "A computer network is a collection of interconnected computing devices that can exchange data and share resources. In a packet-based network, such as the Internet, the computing devices communicate data by dividing the data into small blocks called packets, which are individually routed across the network from a source device to a destination device. The destination device extracts the data from the packets and assemblies the data into its original form. Dividing the data into packets enables the source device to resend only those individual packets that may be lost during transmission.\nCertain network devices, such as routers, maintain tables of routing information that describe routes through the network. A “route” can generally be defined as a path between two locations on the network. Upon receiving an incoming data packet, the network device examines destination information within the packet to identify the destination for the packet. Based on the destination, the network device forwards the packet in accordance with the routing table.\nA network device such as a router uses interface cards (IFCs) for receiving and sending data packets via network links. These IFCs are installed in ports known as interfaces and are configured using interface configurations. A network device can be manufactured with one or more interface configurations, including, for example, ATM, fast Ethernet, and Sonet interface configurations. In particular, a network device may include interface configuration for interfaces in which IFCs have been installed, as well as interfaces in which IFCs have not been installed. In addition, the network device may include a number of different configurations for a particular interface, enabling the network device to be preconfigured with a number of different upgrade paths.\nWhen a new interface card is added to an interface in the network device, the interface is typically configured by selecting one of the pre-defined configurations that is appropriate for the type of interface card added to the interface. For example, if a fast Ethernet interface card is installed, a technician may select a pre-defined fast Ethernet configuration. Because this configuration already exists in the network device, it is not necessary for the technician to load the configuration, and installation is less complicated and time-consuming.\nBefore an interface card is installed in an interface, the interface can be considered a “phantom” interface in that it contains no interface card and is, as a result, inoperative. An interface may also be phantom after removal of a previously installed interface card whose configuration may remain loaded in the network device for potential use in a future installation of an interface card. An interface can also be phantom even if an interface card is installed, for example, if the network device stores multiple configuration options for the pertinent interface card slot. In particular, when the technician selects one of the configuration options, the other configuration options may remain stored in the network device. Generally, an interface can be considered phantom if it is associated with one or more interface card configurations that are not currently active.\nThe presence of a phantom interface may be confusing for installation technicians, network administrators, and other users who access the network device to obtain operational and configuration information. When such a user enters a command to display the configuration of the network device, the network device typically displays all of the configurations currently loaded in the network device, including those for phantom interfaces. This manner of presentation may confuse the user, particularly when multiple configuration options are loaded for a single interface. For example, when an interface card is installed, the user may have difficulty determining which configuration the technician selected for the interface card. Confusion may also arise if no interface card is physically installed in the interface. In particular, a user viewing the configuration may not be aware of the absence of an interface card without physically inspecting the network device interface."}
{"text": "1. Field of the Invention\nThe subject matter disclosed generally relates to the field of robotics used in the medical field.\n2. Background Information\nThere is a growing need to provide remote health care to patients that have a variety of ailments ranging from Alzheimers to stress disorders. To minimize costs it is desirable to provide home care for such patients. Home care typically requires a periodic visit by a health care provider such as a nurse or some type of assistant. Due to financial and/or staffing issues the health care provider may not be there when the patient needs some type of assistance. Additionally, existing staff must be continuously trained, which can create a burden on training personnel. It would be desirable to provide a system that would allow a health care provider to remotely care for a patient without being physically present.\nRobots have been used in a variety of applications ranging from remote control of hazardous material to assisting in the performance of surgery. For example, U.S. Pat. No. 5,762,458 issued to Wang et al. discloses a system that allows a surgeon to perform minimally invasive medical procedures through the use of robotically controlled instruments. There have also been developed “toy” robots for home use. Such robots typically have a relatively simple movement platform and some type of speech synthesis for generating words and sounds. It would be desirable to provide a robotic system that would allow for remote patient monitoring and assistance."}
{"text": "The present invention relates to a dashboard support, and to a method for the production of a dashboard support.\nThe following discussion of related art is provided to assist the reader in understanding the advantages of the invention, and is not to be construed as an admission that this related art is prior art to this invention.\nModern motor vehicles have many driver assist systems as well as comfort and safety systems which are predominantly operated from the inside, in particular via the driver or passenger seats of a motor vehicle. Design features whether relating to the outer look or the interior space gain also increasingly on importance. For that reason, instrument panels, arranged in the motor vehicle between the front A-pillars beneath the windshield, are manufactured of plastic materials and sometimes covered with leather or other coating materials so as to receive a pleasing surface. To stiffen the instrument panel and to render it suitable for supporting various components, a dashboard support is arranged between the A-pillars. The dashboard support represents an elongate, rod-like structure having various mounting areas, for example, for steering column, airbags, glove compartment, heads-up display, navigation system and cable harness. Dashboard supports are typically coupled with the engine compartment rear bulk or the transmission tunnel for increase in stiffness.\nMotor vehicle parts are normally manufactured using manipulators or industrial robots to maneuver the various body components. The various parts are hereby assembled by a multi-stage fully automated process, requiring the parts to be aligned relative to one another for subsequent joining by welding or bonding for example. The body components have openings for engagement by grippers of the manipulators or industrial robots for handling purposes. Dashboard supports are normally handled in the area of the mounting areas for movement and installation in the region of the A-pillars.\nIt would be desirable and advantageous to provide an improved method for the production of a dashboard support to obviate prior art shortcomings."}
{"text": "1. Technical Field\nThe present invention relates to a display device.\n2. Description of the Related Art\nThere have recently been an increasing demand for display devices designed for mobile apparatuses and the like, such as mobile phones and electronic paper. In such display devices, each pixel includes a plurality of sub-pixels that output light of respective colors. The display devices switch on and off the display in the sub-pixels, thereby causing each pixel to display various colors. Display characteristics, such as resolution and luminance, of the display devices have been improved year by year. An increase in the resolution, however, reduces the aperture ratio. To achieve high luminance, it is necessary to increase the luminance of a backlight, resulting in increased power consumption of the backlight. To address this, there has been developed a technology of adding a white pixel serving as the fourth sub-pixel to the conventional red, green, and blue sub-pixels (e.g., Japanese Patent Application Laid-open Publication No. 2010-33014). Because the white pixel increases the luminance, the technology reduces the current value of the backlight, thereby reducing the power consumption.\nJapanese Patent Application Laid-open Publication No. 2010-44389 (JP-A-2010-44389) discloses a light source local dimming control method for controlling dimming of a light source module including a light source block provided with a light source that supplies light to a plurality of image regions. In the method, duty ratios of a first light source and a second light source are primarily determined using a first target luminance value of a first image region closest to a first light source and a second target luminance value of a second image region closest to a second light source adjacent to the first light source. The primarily determined duty ratios are compensated using a target luminance value of a remaining image region excluding the first and the second image regions out of the image regions that receive the light from the first and the second light sources. The first and the second light sources are driven by drive signals resulting from compensation of the primarily determined duty ratios.\nLet us assume a case where the technology disclosed in JP-A-2010-44389 is applied to a sidelight light source including a plurality of light sources at a position facing an incident surface corresponding to at least one side surface of a light guide plate. In this case, the luminance distribution of a backlight may possibly vary in a complicated manner, resulting in unnecessary power consumption.\nFor the foregoing reasons, there is a need for a display device that controls the luminance of light sources of a sidelight light source individually, thereby reducing power consumption of the light sources."}
{"text": "Gas turbine engines typically include a compressor section to pressurize inflowing air, a combustor section to burn a fuel in the presence of the pressurized air, and a turbine section to extract energy from the resulting combustion gases. The compressor section may include a plurality of rotor blades separated by a plurality of stator vane assemblies mounted within the engine casing. Each stator vane assembly may comprise one or more stator shrouds coupled together end to end to form an annular structure. Typically, a small, flat metal part (sometimes referred to as a feather seal) may be inserted horizontally between the connected ends of stator shrouds to minimize air flow leakage from the pressurized gas path. However, gaps on the radially inward and/or radially outward side of the feather seal may still exist as the adjacent stator shrouds separate due to thermal and gas loads, resulting in air flow leakage in the radial and axial direction. Air flow leakage may result in overall loss in performance to the gas turbine engine."}
{"text": "This invention relates generally as indicated to an end closure for composite pressure vessels which provides for the effective transfer of reaction loads between the end closure and composite side wall structure.\nPressure vessels such as used, for example, in fluid actuators for flight controls for aircraft and other high pressure applications may be constructed of composite fiber materials to reduce the weight of the vessels without sacrificing strength. In such pressure vessels, provision must be made for transferring loads that are introduced at the end closures therefor to the load reacting longitudinal fibers of the composite side wall structure. Heretofore, this has been accomplished in various ways such as by adhesively connecting the longitudinal fibers to the end closures, by providing pins on the end closures about which the longitudinal fibers are wrapped, or by providing an interfering ramped surface on the end closures engaged by the ends of the longitudinal fibers and interspersing hoop windings within and around the longitudinal fibers radially outwardly of the ramped surface. Each of these methods may cause difficulties due to the relative inefficiency of the resulting joints. The use of a ramped surface is relatively efficient but does not allow for the transfer of torsional loads from the end closure to the load reacting longitudinal tension fibers."}
{"text": "1. Field of the Invention\nThe present invention relates to a hologram recording medium where a hologram is recorded by overlapping illumination with a recording beam and a reference beam.\n2. Description of the Related Art\nAn example of the conventional hologram recording medium is disclosed in JP-A-H09-305978. The hologram recording medium described in this reference has a configuration in which a transparent dielectric layer is laminated over the entire one surface of a disk-like substrate and a hologram recording layer is laminated thereon. The hologram recording medium as a whole has a disk-like shape. Emboss pits for the track servo are formed at a ratio of 3 pits per 1 track on the boundary surface of the transparent dielectric layer serving as a boundary with the substrate. One emboss pit is positioned exactly on the track, whereas the remaining two emboss pits are located in the positions offset in the opposite directions (radial direction) with respect to the track so as to sandwich the track therebetween. For example, during recording, the boundary surface of the transparent dielectric layer where the emboss pits have been formed is illuminated from the substrate side with a first illumination beam for recording (reference beam) so that the beam is converged on the boundary surface, and a second illumination beam for recording (recording beam) is illuminated from the side of the hologram recording layer so as to overlap the first illumination beam for recording, while being converged in the location on this side of the hologram recording layer. As a result, an interference pattern generated due to overlapping of the first illumination beam for recording and second illumination beam for recording is recorded as a hologram on the hologram recording layer. At this time, the reflected light created by the first illumination beam for recording reflected at the boundary surface of the transparent dielectric layer is converted with a photodetector into a reproduction signal, the optical system is displaced based on this reproduction signal so that the amplitudes corresponding to the two aforementioned offset emboss pits become equal to each other and the track servo is thereby conducted.\nHowever, in the above-described conventional hologram recording medium, the offset quantity of the emboss pits for the track servo is comparatively small. Therefore, the first illumination beam for recording undergoes diffraction in the gap between the emboss pits, and the diffracted beam generated thereby reaches the hologram material layer. As a result, the unnecessary interference pattern created by the diffracted beam is also recorded in the hologram material layer. Another problem is that the first illumination beam for recording serving as the reference beam is intensity modulated by the emboss pits and, therefore, the SN ratio is at a low level during recording and reproduction.\nWith the foregoing in view, it is an object of the present invention to provide a hologram recording medium having a structure in which the recording beam and reference beam are guided in the prescribed direction and making it possible to realize a high SN ratio."}
{"text": "In recent years, portable information devices including touch-panel liquid crystal displays and no physical keyboards, such as tablet PCs and smartphones, have been rapidly widespread. Displays of portable information devices of this type are preferably large in use, but are preferably small size when being carried. In view of this, in a proposed portable information device, a display as well as a chassis is made foldable by using a flexible display such as an organic electro luminescence (EL) display (see, for example, Patent Document 1).\n[Patent Document 1] Japanese Patent Application Laid-Open No. 2014-216025"}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method of recording an image, and more particularly to an image recording apparatus and a method of recording an image which an image formed on a heat-development type photosensitive material by exposure is subjected to heat-development and is then transferred to an image receiving material so as to record the image thereon.\n2. Description of the Related Art\nIn a conventional image recording apparatus, an image formed on a heat-development type photosensitive material is subjected to heat-development and is then transferred to an image receiving material so as to record the image thereon. The image recording apparatus is provided with a photosensitive material magazine for storing a rolled heat-development type photosensitive material. The photosensitive material is pulled out from the photosensitive material magazine by a predetermined length necessary for forming an image thereon by exposure. The photosensitive material thus pulled out is cut and is then conveyed to an exposure section where an image is formed on the photosensitive material by exposure. After completion of the exposure, the photosensitive material is conveyed to a heat-development/transfer section via a water applying section. The image recording apparatus is also provided with an image receiving material magazine for storing a rolled image receiving material. The image receiving material is pulled out from the image receiving material magazine and is then conveyed to the heat-development/transfer section. In the heat-development/transfer section, the photosensitive material and image receiving material are conveyed around a cylindrical heating drum in a state in which they closely contact each other. The photosensitive material and image receiving material conveyed around the heating drum are both heated by heat from a halogen lamp built into the heating drum. As a result, the image formed on the photosensitive material by exposure is heat-developed and is transferred to the image receiving material.\nIn the heat-development section of the image recording apparatus, the photosensitive material and the image receiving material must be conveyed around the cylindrical heating drum. Therefore, a large space is needed to convey the photosensitive material and the image receiving material around the heating drum. Also, since the heating drum has a cylindrical shape, wasted space exists around the heating drum.\nAlso, the image receiving material magazine of the image recording apparatus must have a cylindrical space therein to store a rolled image receiving material. Accordingly, a large space is needed and wasted space exists therein.\nIn the image recording apparatus, the photosensitive material is pulled out from the photosensitive material magazine by a predetermined length necessary for forming an image thereon by exposure, cut, and then conveyed to the exposure section in which an image is formed on the photosensitive material by exposure. Consequently, a large space is necessary to pull out the photosensitive material from the magazine by a predetermined length necessary for forming an image thereon by exposure.\nFor the above-described reasons, the conventional image recording apparatus has the drawbacks of having a large size and high costs."}
{"text": "The construction industry is making increased use of prefabricated trusses in constructing a variety of buildings. Steel buildings have been developed using aircraft structures techniques which may be quickly erected while producing high strength for low weight. Large portions of these steel buildings are fabricated for quick erection, but there is little or no use of prefabricated collapsible steel trusses for these building, or in fact, little use of such trusses in buildings.\nThe major drawbacks to the use of such trusses appears to be size when the trusses are not fully collapsible, difficulty in erection when rigidity is lacking and the ratio of strength to weight of the erected truss. If, however, these problems could be solved, one serious dilemma in the erection of steel and other prefabricated structures would be alleviated.\nAs many erections have painfully learned, high winds during building erection can play havoc with aircraft-type structures which use their metal structural covering for much of their strength. Prefabricated truss members could be preassembled in an essentially flat shape at ground level before being erected in support of a wall or roof. Exposure to potentially disastrous winds would be minimized and greater strength to weight provided without significant increase in erection time.\nIn addition to resolution of this construction industry problem, the light weight compactness of the easily-assembled subject structure lends itself to space applications for such things as manned stations assembled in orbit. The expanding field of oceanography will soon require such structures in marine applications. In the area of portable structures for the military or for campers, the subject invention can provide a shelter from a variety of harsh environments and may be easily erected despite such environmental conditions.\nAccordingly, it is an object of the invention to provide a truss structure which may be collapsed to as compact a form as possible for shipment to an erection site.\nIt is a further object of the invention to provide a truss structure which may be easily and quickly erected from its collapsed form to its truss form.\nIt is an object of the invention to provide a collapsible truss member which when assembled provides rigidity and high strength for low weight.\nIt is also an object of the invention to provide a rod and joint truss structure wherein said joints are simple and inexpensively produced.\nIt is a further object of the invention to produce a collapsible truss which may be erected to form arched, vaulted or compound structural members.\nIt is a further object of the invention to provide a collapsible truss member to be assembled and erected as a part of a larger structure and which will provide rigidity after its erection while other parts of the larger structure are installed."}
{"text": "1. Field of the Invention\nThis invention relates to molding machines, more particularly, to ejector mechanisms commonly used in molding machines. The invention is directed particularly to electrically operated ejector devices having apparatus for selectably increased ejection force.\n2. Description of Related Art\nIn molding machines, plastically deformable material to be molded is formed in cavities defined by mating mold sections and allowed to cure to a state wherein the material will not unacceptably deform upon removal from the mold cavity. The cured material defines molded articles that are removed from the machine upon separation of the mating mold sections. However, as it is common that articles will adhere to one of the mold sections, it is typical to provide ejector pins communicating with the mold cavity and linked to movable ejector members in the mold assembly comprising the mating mold sections. Motion of the ejector pins is effective to dislodge molded articles from the mold section, assuring their complete removal. The movable members are typically translatable and include links to the ejector pins to move the ejector pins between retracted positions whereat the free ends of the ejector pins are flush with mold cavity surfaces and forward positions whereat the free ends protrude into the mold cavity.\nIt is known to use electrically operated actuators to operate ejection devices of molding machines. In particular, it is known to use rotating electrical machines with motion converters to effect linear motion of ejector members. It is known that ejection of articles that have relatively high ratios of internal depth to overall cross sectional area or surface features such as threads can require unusually high ejection force to separate molded articles from so-called mold cores. However, to increase the maximum force exerted by an electrically driven ejection device to accommodate ejection of such articles can result in substantial increases of cost, either as a result of sizing of electrical components and motion converters or to interpose force multipliers between the motor and ejector members. In light of such limitations, there is a need to enable an increase of force exerted on ejector members without the prohibitive attendant costs of known constructions."}
{"text": "The present invention relates to a battery, and more particularly to a battery which has means which, upon the occurrence of an abnormal phenomenon, such as an overcharge or an external short circuit, rapidly operates a safety mechanism, such as a current breaking valve, to ensure the safety of the battery.\nThe spread of portable equipment, such as portable telephones, potable information terminal equipment, and personal computers, has led to an increasing demand for batteries having high energy density. This in turn has led to studies on an increase in capacity. On the other hand, with an increase in the energy density of batteries, an improvement in safety has become important. This is for ensuring safety upon the occurrence of abnormal phenomenon derived from wrong use by a user or malfunction of a circuit. For example, a mechanism (current breaking valve) for breaking current in response to a rise in pressure within the battery and a mechanism (rupture) for releasing internal pressure have been developed and put to practical use. Upon a mistake of connection or malfunction of a charge control circuit, these mechanisms break the current or release the internal pressure before the flow of large current or an increase in battery voltage above a predetermined value causes a dangerous state.\nThe placement of a material capable of evolving a gas within the battery has been proposed to rapidly increase the pressure within the battery upon the occurrence of an abnormal phenomenon. For example, in a rechargeable battery with a nonaqueous electrolysis solution, the addition of lithium carbonate to a positive electrode has been proposed to enable lithium carbonate to be decomposed upon overcharge to evolve carbon dioxide gas which operates the current breaking valve.\nLithium carbonate as the conventional material, however, is decomposed upon an increase in voltage as a result of the overcharge, and cannot respond to the occurrence of an abnormal phenomenon, which lowers voltage, for example, external short circuit. Further, the gas evolved is carbon dioxide which is easily dissolved in the electrolysis solution. Therefore, the pressure rise rate upon the occurrence of abnormal phenomenon is not necessarily satisfactory.\nIn view of the above problems of the prior art, the present invention has been made, and it is an object of the present invention to provide a battery which has means which, upon the occurrence of an abnormal phenomenon, such as an overcharge or an external short circuit, rapidly operates a safety mechanism, such as a current breaking valve, to ensure the safety of the battery and which, particularly upon the occurrence of an abnormal phenomenon, can surely and stably increase the pressure within the battery.\nThe above object of the present invention is attained by a battery comprising a positive electrode, a negative electrode, a separator, an electrolysis solution, and a hermetically sealed container,\nthe hermetically sealed container containing in its interior a compound represented by formula (1):\nXxe2x80x94Oxe2x80x94COxe2x80x94Rxe2x80x83xe2x80x83(1)\nwherein X represents a group which, upon decomposition of the compound caused by a rise in temperature, is eliminated to evolve a gas insoluble or slightly soluble in the electrolysis solution; and R represents a group which controls the decomposition temperature of the compound, the group R having, in the form of a compound represented by chemical formula RH1, a hydrogen atom H1 with an acid dissociation equilibrium constant (pKa) of not more than 13.\nAccording to the present invention, the group X in formula (1) preferably represents a substituted or unsubstituted alkyl group having 3 to 7 carbon atoms and is specifically selected from the group consisting of substituted or unsubstituted propyl, butyl, and pentyl groups.\nAccording to a preferred embodiment of the present invention, the hydrogen atom H1 in the compound represented by chemical formula RH1 is bonded to an oxygen or nitrogen atom.\nAccording to the present invention, the decomposition temperature of the compound represented by formula (1) can be regulated in the range of 100 to 150xc2x0 C. by properly selecting the group R.\nFor example, in formula (1), when the group X represents a tert-butyl group, the compound is rapidly decomposed upon a temperature rise to evolve carbon dioxide, 2-methylpropene, and a compound represented by chemical formula RH1. In this case, since 2-methylpropene has low solubility in both water electrolysis solution and nonaqueous electrolysis solution, the pressure within the battery can be rapidly increased to operate the current breaking valve or the rupture."}
{"text": "The present invention relates to a continuously variable transmission means for use in a vehicle. The continuously variable transmission (CVT) varies transmitting speed ratios of a revolution number of a driven pulley to that of a driving pulley. More particularly, the present invention relates to a transmission means having a control means which permits a low engine RPM, even when an engine experiences a large load.\nFIG. 1 shows a graph illustrating a relationship between an engine torque and an engine speed (RPM) of a vehicle in which a continuously variable transmission means is mounted. Heretofore, the RPM of an engine equipped with a CVT was designed such that it assumed a value determined by the curves \"a\" and \"b\" in FIG. 1. It is preferable that an engine RPM is determined so that the amount of fuel consumption be minimized over the range. The entire curve \"a\" shown in FIG. 1 is determined such that the amount of fuel consumed may be minimized. On the other hand, in general, vibrations in the engine increase when the engine is at a at a low RPM, while experiencing a high load. This results in the occurrence of undesirable vibration in vehicle body, and noise in a passenger compartment. Further, this vibration shortens the endurance of driving parts employed in a vehicle. It is quite difficult to determine an engine RPM especially under the condition that an engine RPM is low and an engine load is high. Hence, an engine RPM must be varied from an engine speed curve, indicated by the reference \"a\", to an another engine speed curve indicated by a reference \"b\" during the time when engine RPM is low. A point D1 on the curve \"a\" is a point where the curve \"a\" intersects with the curve \"b\". The curve \"b\" is determined independently from the minimum fuel consumption curve \"a\", as shown in FIG. 1. The preferable engine characteristics at a low engine RPM under a high engine load varies, according to the type of engine mounting structure to a vehicle body, the vehicle body structure and the engine operating conditions. Hence, in order to obviate the above-mentioned drawbacks the engine characteristic curve at a low engine speed is determined by the curve \"b\", which is located at relatively higher engine RPM's than that of the curve \"c\". This results in a large amount of a gasoline consumption and impairs efficient operation."}
{"text": "1. Technical Field\nThe present invention relates in general to electronic communications and, in particular, to controlling throughput of message requests in a messaging system. Still more particularly, the present invention relates to specifying output of a message request according to a priority requirement set by a user receiving the message request.\n2. Description of the Related Art\nAs the Internet and telephony expand, the ease of communications between individuals in different locations continues to expand as well. One type of electronic communication is supported by messaging which includes the use of computer systems and data communication equipment to convey messages from one person to another, as by e-mail, voice mail, unified messaging, instant messaging, or fax.\nWhile e-mail has already expanded into nearly every facet of the business world, other types of messaging continue to forge into use. For example, instant messaging systems are typically utilized in the context of an Internet-supported application that transfers text between multiple Internet users in real time.\nIn particular, the Internet Relay Chat (IRC) service is one example of instant messaging that enables an Internet user to participate in an on-line conversation in real time with other users. An IRC channel, maintained by an IRC server, transmits the text typed by each user who has joined the channel to the other users who have joined the channel. An IRC client shows the names of the currently active channels, enables the user to join a channel, and then displays the other channel participant's words on individual lines so that the user can respond.\nSimilar to IRC, chat rooms are often available through on-line services and provide a data communication channel that links computers and permits users to converse by sending text messages to one another in real-time.\nSome messaging systems allow users sending messages to assign a priority level, such a high, medium, or low, to a message. However, a limitation of such systems is that the priority level is typically only displayed with the message when received by another user utilizing the same software utilized to create the message. Further, a limitation of such systems is that the message prioritization is set only by the user sending the message. Further, the user is still typically interrupted either by a visual or audible cue whether the message is low priority or high priority.\nAs instant messaging systems and telephony based messaging systems expand, the number of unwanted and often unnecessary daily intrusions increases. Other users are enabled to see when another user has logged on and send instant messaging requests that take up screen space and are both audibly and visibly intrusive.\nSome messaging systems enable a user to select from among several “away messages” that are transmitted automatically in response to a another user's request to talk. With an “away message” set, a user may continue working while logged on, however other users are informed that the user is not available. With telephony devices, a user is provided with an answering service, rather than having to answer a call.\nHowever, messaging systems and telephony devices are limited in that while a particular user may indicate to other requesting users that the particular user is unavailable, in order for that user to utilize a network connection, the user is still intruded upon by the requests to talk. For example, in an instant messaging system, an automatic “not available” is transmitted, but either all requests still come through or no requests come through. Similarly, a ringer may be turned off on a telephony device such that the caller ID is displayed, but then a user is not alerted to any calls.\nA further disadvantage of “away messages” is that a single message is displayed to all users requesting to talk. However, it would be advantageous to tailor “away messages” depending on the person receiving the message such that a spouse receives one message while a boss receives another.\nIn addition, whether a user is designated as away or not, most messaging systems only offer one format for notifying users of a request to talk. Such requests often utilize an unsatisfactory amount of screen space and may include unwanted audible prompts.\nIn view of the foregoing, it would be advantageous to provide a method, system, and program for controlling the throughput of message requests through a messaging system such that a user may specify which message requests are output and how the message request is output. Further, it would be advantageous to provide a method, system, and program for automatically specifying the response message transmitted to the user requesting to chat."}
{"text": "The bandwidth of content being contributed to network users by service providers of metro and core networks is increasing rapidly. A main factor in this increase is new video services, especially High-Definition TV (HDTV). With modern compression techniques (such as MPEG-4, H.264 and VC-1), HDTV can be transmitted at 8-9 Mbps, and a Standard-Definition TV (SDTV) channel at 2-1.5 Mbps.\nThe metro and core networks that provide connections between the head ends and the central offices are typically based on optical fiber infrastructures, so they are usually able to cope with this increase in bandwidth. The end-user access network, however, is typically based on twisted-pair copper wires over which Digital Subscriber Line (xDSL) technology is used to provide communication rates of up to 10 Mbps with Asymmetric Digital Subscriber Line 2 Plus (ADSL2+, also known as International Telecommunication Union standard G.992.5), which allows for a single HDTV channel.\nThe network is represented in hierarchical fashion, with the following layers:                Core Layer        Aggregation Layer (also referred to as the “Metro Layer”)        Access Layer        \nThe end-user is in the Access Layer. The term “access edge” herein denotes the boundary between the Aggregation Layer and the Access Layer of a network. As a non-limiting example, a Digital Subscriber Line Access Multiplexer (DSLAM) with ADSL2 output to a copper twisted pair medium, as described above, exists in the Access layer on the access edge. Other devices besides a DSLAM may be used to accomplish the processing necessary for delivering data to end-users; in the descriptions and drawings herein, a DSLAM is used only as a non-limiting example of such equipment, for purposes of illustration, and is not intended to restrict the present invention to any specific configuration or use.\nCurrently, the average number of television sets per household is between 2 and 3, often requiring multiple different HDTV channels simultaneously. As a result, video over xDSL may not have sufficient bandwidth to compete with cable and satellite offerings.\nOne solution is to bring optical fiber into the end-user's premises (or very close to it), but the capital investment required for such a solution is still very high."}
{"text": "The 3rd Generation Partnership Project (3GPP) long term evolution (LTE) system, introduced as 3GPP release 8, is an improved universal mobile telecommunication system (UMTS). An LTE system offers high peak data rates, low latency, improved system capacity, and low operating cost resulting from simple network architecture. In the LTE system, an evolved universal terrestrial radio access network (E-UTRAN) includes a plurality of evolved Node-Bs (eNBs) communicating with a plurality of mobile stations, referred as user equipment (UE). 3GPP introduces new features to help LTE system operators to further optimize network planning in a cost-effective way. Minimization of Drive Test (MDT) is one of the features where UEs collect measurements and report measurement information to their serving eNBs.\nMDT has been worked on in 3GPP to help with network optimization. Network optimization is traditionally done by manual drive testing, which is costly and causes additional CO2 emissions. The essence of MDT feature is to provide for normal mobile terminals the possibility to record and log information relevant to the radio communication of the mobile terminals, correlated with the geographical location of the mobile terminals. MDT feature enables UEs to perform Operations, Administration, and Maintenance (OAM) activities, such as neighborhood detection, measurements, logging and recording for OAM purposes, which includes radio resource management (RRM) and optimization purposes.\nWhile studying the possible scope of minimization of drive test (e.g., 3GPP TR 36.805), it has been found beneficial to include Quality of Service (QoS) measurements for MDT, and in particular a throughput measurement. It is assumed that it would be desirable to as much as possible mimic manual drive testing. It would be desirable to provide the possibility to measure throughput, and while the UE pass through a network doing frequent handovers. Note that it is expected that frequent handover may affect TCP performance and lead to reduced user-experienced throughput.\nSeveral definitions of throughput measurements exist in 3GPP TR36.314, TR 32.425, and TR 32.451. However, there is no definition supports well all the MDT requirements. It is an objective of the current invention to propose a throughput measurement definition, which fulfills well all the MDT requirements with minimum complexity."}
{"text": "The United States government may have rights in this invention pursuant to funding arrangements with the Department of Defense.\n1. Field of the Invention\nThis invention relates to an improvement in electromagnetic railgun construction which substantially improves railgun performance and permits use of a solid armature. In particular, it relates to a railgun having a layer of low conductivity material, such as graphite, along the railgun bore which inhibits the concentration of current density along the interface of the rail and armature.\n2. Description of the Related Art\nThermodynamic guns are widely used and generally understood in a broad context. In an ordinary thermodynamic gun, a propellant burns to generate high pressure gas that pushes a projectile down a bore. While thermodynamic guns are used in many applications besides weapons --for example scientific and industrial applications-- their use is somewhat limited because of the maximum velocities attainable. Thus, physical limitations limit the projectile from such thermodynamic guns from reaching velocities much greater than one kilometer per second.\nElectromagnetic railguns have been widely investigated since World War II as an alternative to thermodynamic guns because of the possibilities of achieving projectile hypervelocities (greater than one kilometer per second).\nThe early electromagnetic railguns incorporated a solid armature which was propelled between the rails by the electromagnetic force generated by the current flow through the armature and the rails. However, it was soon found that at high speeds around one kilometer per second, the rails and armature were substantially damaged, possibly as a result of ohmic heating and/or internal forces. Further, increases in current flow tended to only increase rail and armature gouging without an increase in armature velocity. Thus, armature velocities in excess of one kilometer per second were not practically attainable for railguns using solid armatures.\nIn the early 1970's, R. A. Marshall, J. P. Barber, and others at the Australian National University, Canberra, Australia, developed railguns using plasma armatures which could obtain hypervelocities and could make efficient use of high current, pulsed power supplies, such as compulsators and homopolar generators. Plasma armature railguns are generally described in S. C. Rashleigh and R. A. Marshall, \"Electromagnetic Acceleration of Macroparticles to High Velocities,\" 49 J. App. Phys. 2540 (Apr. 1978) (incorporated herein by reference). Such power supplies are broadly illustrated in U.S. Pat. Nos. 4,200,831; 4,459,504; 4,246,507; and U.S. Pat. Application Ser. No. 06/689,868 (incorporated herein by reference).\nIn recent years, additional research has revealed numerous problems associated with very high current plasma armatures. For example, at the high currents necessary to obtain hypervelocities, rail errosion has been a significant problem which essentially relegates the railgun to a one or two shot application. Additionally, plasma armature type railguns require a sealed bore (open only at the muzzle) capable of withstanding the substantial electromagnetic forces generated; the gaskets, seals, and insulator material associated with such bores have proven a significant problem. For example, in addition to rail and insulator damage, metallic deposits often adhere to the insulator surfaces after firing, causing arcing problems during subsequent firing.\nIn addition to the practical difficulties of plasma-type electromagnetic railguns, several fundamental problems have become of increasing concern. For example, at the typical working currents in question, the dissipative armature voltage drop is an order of magnitude greater than desirable. Perhaps more fundamentally, plasma armatures are designed to accelerate a projectile using base pressures. Base pressure acceleration (such as also used in thermodynamic guns) places severe design limits on the projectile. For example, a projectile must be able to withstand the extreme temperature and pressures exerted at its base by the driving plasma. Such projectile design limits are removed and gun efficiency substantially improved if a solid armature is used as the projectile. A solid armature/projectile offers the prospect of acceleration by body forces on the armature/projectile.\nWhile a solid armature railgun avoids many of the problems associated with plasma armatures--e.g. substantial dissipative armature voltage drop and base pressure acceleration--attempts to obtain hypervelocities with solid armature railguns have proven unsuccessful. Thus, it would a significant advance in the art if a railgun were devised which could utilize a solid armature and operate at velocities substantially in excess of one kilometer per second."}
{"text": "This invention relates to a new and distinct cultivar of the Saxifragaceae family. The botanical name of the plant is Hydrangea macrophylla (Thunb.) ‘Galilee’.\nThe new cultivar originated as a seedling from a controlled cross between Hydrangea macrophylla (Thunb.) True Blue - U.S. Plant Pat. No. 18,593 (the seed parent), and the unpatented commercial variety Hydrangea macrophylla (Thunb.) ‘Mathilda Gutges’ (the pollen parent).\nThe variety ‘Galilee’ has stronger stems than ‘Mathilda Gutges’ and larger florets than either ‘True Blue’ or ‘Mathilda Gutges’. The sepals of the new variety develop a deeper blue tone with the addition of less aluminum than its parent ‘True Blue’.\nTABLE 1U.S. Plant Pat.UnpatentedNew VarietyNo. 18,593‘Mathilda‘Galilee’‘True Blue’Gutges’Stem strengthStrongStrongStrongPigmentationR.H.S. 100BR.H.S 100DR.H.S. 83Cof Sepals(blue group)(blue group)(violet group)Easier toat edges ofproduce a moresepals.blueR.H.S. 98Apigmentation(violet-bluewith lessgroup) inalumina thancenter ofparentssepals.Sepal marginsDentateDentateSerrateFloret SizeLarger thanSmaller thanSmaller thanParents‘Galilee’‘Galilee’\nAsexual reproduction was first accomplished when vegetative cuttings were taken from the initially selected plant. Examination of asexually reproduced, successive generations grown at a nursery at Half Moon Bay, Calif., under the direction of the inventor show that the combination of characteristics as herein disclosed for ‘Galilee’ remains firmly fixed through three generations."}
{"text": "1. Technical Field\nThe present disclosure is directed to a dynamic stabilization device and system for spinal implantation and, more particularly, to a dynamic stabilization device and system that is adapted to be positioned/mounted relative to first and second laterally-spaced pedicle screws and that includes at least one dynamic stabilization member that is positioned beyond the region defined between the pedicle screws, e.g., in an “overhanging” orientation.\n2. Background Art\nLow back pain is one of the most expensive diseases afflicting industrialized societies. With the exception of the common cold, it accounts for more doctor visits than any other ailment. The spectrum of low back pain is wide, ranging from periods of intense disabling pain which resolve, to varying degrees of chronic pain. The conservative treatments available for lower back pain include: cold packs, physical therapy, narcotics, steroids and chiropractic maneuvers. Once a patient has exhausted all conservative therapy, the surgical options range from micro discectomy, a relatively minor procedure to relieve pressure on the nerve root and spinal cord, to fusion, which takes away spinal motion at the level of pain.\nEach year, over 200,000 patients undergo lumbar fusion surgery in the United States. While fusion is effective about seventy percent of the time, there are consequences even to these successful procedures, including a reduced range of motion and an increased load transfer to adjacent levels of the spine, which accelerates degeneration at those levels. Further, a significant number of back-pain patients, estimated to exceed seven million in the U.S., simply endure chronic low-back pain, rather than risk procedures that may not be appropriate or effective in alleviating their symptoms.\nNew treatment modalities, collectively called motion preservation devices, are currently being developed to address these limitations. Some promising therapies are in the form of nucleus, disc or facet replacements. Other motion preservation devices provide dynamic internal stabilization of the injured and/or degenerated spine, without removing any spinal tissues. A major goal of this concept is the stabilization of the spine to prevent pain while preserving near normal spinal function. The primary difference in the two types of motion preservation devices is that replacement devices are utilized with the goal of replacing degenerated anatomical structures which facilitates motion while dynamic internal stabilization devices are utilized with the goal of stabilizing and controlling abnormal spinal motion.\nOver ten years ago a hypothesis of low back pain was presented in which the spinal system was conceptualized as consisting of the spinal column (vertebrae, discs and ligaments), the muscles surrounding the spinal column, and a neuromuscular control unit which helps stabilize the spine during various activities of daily living. Panjabi M M. “The stabilizing system of the spine. Part I. Function, dysfunction, adaptation, and enhancement.” J Spinal Disord 5 (4): 383-389, 1992a. A corollary of this hypothesis was that strong spinal muscles are needed when a spine is injured or degenerated. This was especially true while standing in neutral posture. Panjabi M M. “The stabilizing system of the spine. Part II. Neutral zone and instability hypothesis.” J Spinal Disord 5 (4):390-397, 1992b. In other words, a low-back patient needs to have sufficient well-coordinated muscle forces, strengthening and training the muscles where necessary, so they provide maximum protection while standing in neutral posture.\nDynamic stabilization (non-fusion) devices need certain functionality in order to assist the compromised (injured or degenerated with diminished mechanical integrity) spine of a back patient. Specifically, the devices must provide mechanical assistance to the compromised spine, especially in the neutral zone where it is needed most. The “neutral zone” refers to a region of low spinal stiffness or the toe-region of the Moment-Rotation curve of the spinal segment (see FIG. 1). Panjabi M M, Goel V K, Takata K. 1981 Volvo Award in Biomechanics. “Physiological Strains in Lumbar Spinal Ligaments, an in vitro Biomechanical Study.” Spine 7 (3):192-203, 1982. The neutral zone is commonly defined as the central part of the range of motion around the neutral posture where the soft tissues of the spine and the facet joints provide least resistance to spinal motion. This concept is nicely visualized on a load-displacement or moment-rotation curve of an intact and injured spine as shown in FIG. 1. Notice that the curves are non-linear; that is, the spine mechanical properties change with the amount of angulations and/or rotation. If we consider curves on the positive and negative sides to represent spinal behavior in flexion and extension respectively, then the slope of the curve at each point represents spinal stiffness. As seen in FIG. 1, the neutral zone is the low stiffness region of the range of motion.\nExperiments have shown that after an injury of the spinal column or due to degeneration, neutral zones, as well as ranges of motion, increase (see FIG. 1). However, the neutral zone increases to a greater extent than does the range of motion, when described as a percentage of the corresponding intact values. This implies that the neutral zone is a better measure of spinal injury and instability than the range of motion. Clinical studies have also found that the range of motion increase does not correlate well with low back pain. Therefore, the unstable spine needs to be stabilized especially in the neutral zone. Dynamic internal stabilization devices must be flexible so as to move with the spine, thus allowing the disc, the facet joints, and the ligaments normal physiological motion and loads necessary for maintaining their nutritional well-being. The devices must also accommodate the different physical characteristics of individual patients and anatomies to achieve a desired posture for each individual patient.\nWith the foregoing in mind, those skilled in the art will understand that a need exists for a spinal stabilization device which overcomes the shortcoming of prior art devices. The present invention provides such an apparatus and method for spinal stabilization."}
{"text": "1. Field of Invention\nThis invention relates to running operating system in device(s), including computer system(s) or computer-controlled device(s) or operating-system-controlled device(s). Such device(s) is/are mentioned hereafter as Device(s).\nIn particular, this invention relates to operating systems to which this invention can be and is applied; this invention being a method, capable of being implemented in computer-executable or operating-system-executable instructions or programmes, providing for running operating system on virtual system memory address space, being dynamically mapped onto a storage medium, such as but not limited to internal physical memory or internal Random Access Memory (RAM or ram used hereafter), in or accessible to the Device(s) capable of running the operating system under concern.\nIn this relation, it makes possible, in the Device(s) capable of running the operating system under concern, the phenomenon of booting up and running operating system with the use of virtual disk driver on virtual system memory, the virtual system memory being dynamically mapped onto storage medium for its access.\n2. Description of Related Arts\nSince Microsoft Windows is the most popular operating system used nowadays, it is taken here as an example for illustrating how this invention can be carried out. This invention however can be applied to other operating systems with similar designs.\nNormally, operating systems run up from and on fixed internal IDE/SCSI hard disk(s) of a local computer system or operating-system-controlled device or on the network fileserver of a network computer. Attempts however have been made to run operating systems in other ways. Using Microsoft Windows as an example, the solution of creating a CD capable of booting the full Microsoft Windows 98 at http://www.ct.heise.de/ct/english/99/! 1/206/ put forward by Tobias Remberg and Hajo Schulz provides a suggested solution for creating a bootable CD for running up Microsoft Windows 98. Besides this, Microsoft Windows Pre-installed Environment now can also run up on a live CD or DVD or be loaded up from hard disk to run as ramdisk-based. So it is now a common phenomenon to boot up Microsoft Windows from a local IDE/SCSI hard disk, from a CD or DVD, from a USB storage device or from a network fileserver or in hand-held device(s) to run as a ramdisk-based operating system.\nSo given time and will, Microsoft Windows can be made to boot up from more and more types of storage medium through designing and implementing the corresponding boot-up device drivers for accessing such types of storage medium. For certain type of storage medium, such as but not limited to internal RAM within a computing system, to be used for storing and booting up the system image of certain versions of Microsoft Windows, such as Microsoft Windows 2000 and XP, boot-up device driver has to be designed and implemented in such a way that the storage medium, for instance the internal RAM, is made available to the operating system for boot-up and subsequent use through allocating virtual memory out of the system memory pool and mapping such virtual memory onto the storage space of the internal RAM."}
{"text": "1. Technical Field\nThe present invention relates to a technology that cleans the nozzle forming surface of a discharge head which discharges liquid, and, in particular, to a technology that performs wiping with a wiping member on the nozzle forming surface using cleaning liquid.\n2. Related Art\nIn the related art, a liquid discharging apparatus, such as an ink jet printer, which discharges liquid, such as ink or the like, from the nozzles of a discharge head is known. In such an apparatus, there is a case in which the liquid discharged from the nozzles adheres to the nozzle forming surface of the discharge head. If the liquid adheres to the nozzle forming surface, there is a problem in that the liquid is not appropriately discharged from the nozzles, thereby causing deterioration of image quality. Here, various technologies for removing liquid which adheres to the nozzle forming surface of a discharge head are proposed.\nFor example, a head cleaning mechanism disclosed in JP-A-2009-233896 is provided with an ejecting unit that ejects cleaning liquid onto a nozzle forming surface (droplet discharge surface) in order to suitably remove liquid which adheres to the nozzle forming surface. Further, after the cleaning liquid is ejected onto the nozzle forming surface, wiping is performed by a wiping member, and thus it is possible to suitably clean the nozzle forming surface.\nAs described above, in the head cleaning apparatus disclosed in JP-A-2009-233896, the cleaning liquid is directly ejected onto the nozzle forming surface, and thus the cleaning liquid easily comes in nozzles. Therefore, for example, there is a problem in that meniscus which are formed in ink in the vicinity of the openings of the nozzles are destroyed by the cleaning liquid which comes in the nozzles, and thus it is difficult to appropriately discharge ink from the nozzles. Here, in JP-A-2009-233896, it is described that the cleaning liquid is ejected onto the region of the nozzle forming surface other than the nozzles, and thus the cleaning liquid is prevented from coming in the nozzles. However, if the cleaning liquid is ejected onto the nozzle forming surface, there is limitation of preventing the cleaning liquid from coming in the nozzles."}
{"text": "As disclosed in Froeschle U.S. Pat. No. 4,456,872, hysteretic current-mode control offers significant advantages for many applications. Hysteretic current-mode control offers the tightest and most accurate control of inductor current, is unconditionally stable regardless of duty cycle, and offers excellent transient response to step loads. It is also inherently load-current limiting. Other advantages include short-circuit-proof operation in buck-derived topologies, easy loop-stability design, absence of subharmonic oscillation, instantaneous peak-to-average inductor-current ratio, and the absence of a need for slope compensation for duty cycles above 50%.\nAlthough hysteretic current-mode control offers many advantages, it has not yet achieved the widespread use which might be anticipated given the advantages it offers. This may be due to the fact that most hysteretic current-mode control circuits of the present state of the art are complex, requiring numerous components, with an accompanying significant cost disadvantage with respect to simpler, non-hysteretic arrangements. A hysteretic current-mode control circuit requiring fewer and simpler components would naturally be less expensive to design and manufacture, and would make the advantages of this mode of control available in more applications.\nThe present invention recognizes the need for a circuit for achieving hysteretic current-mode control which can be manufactured using a small number of simple components to form a low cost unit suitable for use in applications such as a low cost battery charger, a regulated voltage source in a battery operated electronic device or any other application requiring a low cost DC-to-DC converter, as more fully described below."}
{"text": "This application is based upon and claims the benefit of priority from the prior Japanese Patent Application No. 2000-294966, filed on Sep. 27, 2000; the entire contents of which are incorporated herein by reference.\nThe present invention relates to a semiconductor device and a manufacturing method thereof, and more particularly to a structure of a MIS type integrated circuit on an SOI substrate and a manufacturing method thereof.\nIt has become an important target over the recent years to decrease a consumption of the electric power of an LSI, and a MIS type transistor using an SOI (Silicon On Insulator) substrate has been increasingly developed as one of effective solving methods.\nFIGS. 11A through 11G are sectional views showing manufacturing steps of a method of manufacturing the conventional MIS type transistor provided on the SOI substrate.\nTo start with, as shown in FIG. 11A, an SOI substrate is composed of a silicon substrate 1, an embedded oxide layer 2 and a mono-crystalline layer. On this SOI substrate, a device isolation is carried out by a STI (Shallow Trench Isolation) device isolation method based on a normally used shallow trench, and the isolation is done with a isolation region 5 embedded with an oxide layer, thereby providing a device region 3 of which the surface is covered with a thermal oxide layer 4.\nNext, as shown in FIG. 11B, an NMOS channel region 6 is provided by implanting ions of a p-type impurities into the device region 3, and a PMOS channel region 7 is provided by implanting ions of an n-type impurities into the device region 3 by making use of resist patterning based on photolithography.\nSubsequently, as shown in FIG. 11C, after removing the oxide layer 4 on the device regions 6, 7, a gate insulation layer 8 is provided by performing again the thermal oxidation. Thereafter, undoped polysilicon is deposited over the entire surface by use of an LPCVD method, and gate electrodes 9, 10 are provided on an NMOS channel region 6 and a PMOS channel region 7 by using the resist patterning and reactive ion etching (RIE).\nNext, as shown in FIG. 11D, a low concentration diffused layer 11 is provided in the NMOS region by implanting the ions of n-type impurities, and a low concentration diffused layer 12 is provided in the PMOS region by implanting the ions of p-type impurities by using the resist patterning(not shown) and the gate electrodes 9, 10 as ion implantation masks.\nNext, as shown in FIG. 11E, a nitride layer is deposited by using the LPCVD method, and thereafter a side wall 13 is provided on side surfaces of the gate electrodes 9, 10 by effecting the RIE. Further, a high concentration diffused layer 14 is provided in the NMOS region by implanting the ions of the n-type impurities, and a high concentration diffused layer 15 is provided in the PMOS region by implanting the ions of the p-type impurities by making the use of the resist patterning (not shwon). The ion implantation for providing these high concentration diffused layers 14, 15 also functions to add the impurities to the gate electrodes 9, 10 at the same time. That is, the gate electrode 9 becomes an n-type gate electrode, and the gate electrode 10 becomes a p-type gate electrode. Thereafter, a refractory metal such as Co, Ti, Ni is deposited on an entire wafer surface, and the thermal process is executed thereon, thereby providing a metal silicide 16 selectively on only the region where the silicon of the MIS type transistor is exposed.\nSubsequently, as shown in FIG. 11F, after an oxide layer 17 serving as an inter-layer insulating layer has been deposited, a contact hole 18 to the NMOS, a contact hole 19 to PMOS and a contact hole 20 to the a silicon substrate 1 are respectively formed by the resist patterning based on the photolithography and by the RIE. The contact hole is formed deeper by thickness of the embedded oxide layer 2 and thickness of the isolation region 5 than the contact holes 18, 19.\nFurther, a high concentration impurity layer 21 having the same conductivity as that of the silicon substrate 1 is provided at the bottom of the contact hole 20 by effecting the resist patterning and the ion implantation, and thereafter the thermal process executed thereon, thereby activating the impurity layer 21. Note that this contact hole 20 is used for stabilizing an electric potential of the silicon substrate 1, and the high concentration impurity layer 21 is provided to ensure an ohmic contact between a interconnection metal formed subsequently and the silicon substrate 1.\nThereafter, a metal 22 of Ti/TiN etc is deposited thin within the contact holes 18, 19, 20, and a metal 23 such as tungsten (W) etc is grown with the metal 22 used as a base, and polishing is executed, whereby the metals 22, 23 are left only within the contact holes. Further, after the metal interconnections of Al etc has been deposited on the entire wafer surface, a predetermined metal interconnection 24 is formed by effecting the resist patterning and the RIE.\nThus, according to the conventional semiconductor device, the MIS type transistor is provided and the inter-layer insulation layer is deposited. Thereafter, when forming the contact hole with respect to each of the electrodes to the transistor, the contact holes are so opened as to penetrate the inter-layer insulation layer, the device isolation oxide layer and the embedded oxide layer at the same time, and the metal interconnection is provided as in the case of other contacts, thereby making an electric connection to the silicon substrate. Accordingly, in the semiconductor device manufactured by using the present method, the circuit using the MIS type transistor can be provided on the SOI substrate, which largely contributes to decrease the consumption of the electric power of the LSI.\nThere arise, however, the following problems inherent in the conventional semiconductor device and the manufacturing method thereof described above.\nFirst, it is desirable that the metal interconnection for stabilizing the electric potential of the silicon substrate be, when forming this metal interconnection, taken from the underside of the substrate, however, there is no alternative but to often take it out from the surface of the substrate as the case may be. In this case, as shown in FIG. 11F, the contact hole 20 to the silicon substrate is formed deeper by the thickness of the device isolation oxide layer plus the embedded oxide layer than other contact holes 18, 19, and hence an aspect ratio increases, with the result that it is difficult to form a minute hole corresponding to a design rule.\nSecond, even when forming the contact hole, it is quite difficult to grow a base metal thin layer for growing tungsten uniformly up to side and bottom portions of the hole having a high aspect ratio as in the case of the contact hole 20.\nThird, it is required for ensuring the ohmic contact with respect to the silicon substrate that the high concentration impurity layer 21 be provided at the bottom of the contact hole 20. If the impurity layer formed by executing the ion implantation is activated at a temperature as high as, e.g., 950xc2x0 CRTA after the contact hole 20 has been formed, however, the metal silicide already provided on the MIS type transistor is weak to the heat and decline in its characteristic. This results in a problem that a junction leak is brought about, and a resistance rate rises.\nAs described above, in the semiconductor device manufactured by using the conventional technology, it is difficult to stably form the contact hole to the silicon substrate constituting the SOI substrate.\nA semiconductor device according to an embodiment of the present invention comprises an SOI substrate formed with a mono-crystalline semiconductor layer through an embedded insulating layer on a first conductivity type semiconductor substrate; a MIS type field-effect transistor provided within a device region defined by isolating said mono-crystalline semiconductor layer with a device isolation region and having a gate electrode; an opening formed penetrating said device isolation region and said embedded insulating layer and reaching said semiconductor substrate; and a substrate electrode provided covering said opening and taken out up to the surface flush with said gate electrode.\nA semiconductor device according to another embodiment of the present invention comprises: an SOI substrate configured by providing a mono-crystalline semiconductor layer through an embedded insulation layer on a first conductivity type semiconductor substrate of which the surface is provided with a first impurity diffused layer of the first conductivity type and a second impurity diffused layer of the second conductivity type; a MIS field-effect transistor provided within a device region defined by isolating said mono-crystalline semiconductor layer with a device isolation region, and having a gate electrode;\nfirst and second openings penetrating said device isolation region and said embedded insulation layer and reaching said first impurity diffused layer of said semiconductor substrate; third and fourth openings penetrating said device isolation region and said embedded insulation layer and reaching said second impurity diffused layer of said semiconductor substrate; first and fourth polysilicon substrate electrodes of the first conductivity type, taking such a configuration as to fill and cover said first and fourth openings, and taken out up to the surface flush with said gate electrode; and second and third polysilicon substrate electrodes of the second conductivity type, taking such a configuration as to fill and cover said second and third openings, and taken out up to the surface flush with said gate electrode, wherein first and second diodes for an electrostatic protection are composed respectively of a group of said first polysilicon electrode, said first impurity diffused layer and said second polysilicon electrode, and a group of said third polysilicon electrode, said second impurity diffused layer and said fourth polysilicon electrode.\nA method for manufacturing semiconductor device according to an embodiment of the present invention comprises: providing a device isolation region for defining a device region on a mono-crystalline semiconductor layer of an SOI substrate formed with a mono-crystalline semiconductor layer through an embedded insulation payer on a semiconductor substrate of a first conductivity type; forming an opening penetrating said device isolation region and said embedded insulation layer and reaching said semiconductor substrate; depositing polysilicon on said SOI substrate and within said opening and providing a gate electrode and a substrate electrode of said MIS type field-effect transistor by executing the patterning thereon; and implanting impurities into said gate electrode and said substrate electrode.\nA method for manufacturing semiconductor device according to another embodiment of the present invention comprises: providing a device isolation region for defining a device region on a mono-crystalline semiconductor layer of an SOI substrate formed with a mono-crystalline semiconductor layer through an embedded insulation payer on a semiconductor substrate of a first conductivity type; providing a first polysilicon layer on the entire surface; forming an opening penetrating said first polysilicon layer, said device isolation region and said embedded insulation layer and reaching said semiconductor substrate;\nproviding a second polysilicon layer on the entire surface to cover said opening; providing a substrate electrode by patterning said first and second polysilicon layers; and\nimplanting an impurity into said substrate electrode."}
{"text": "Reaction receptacles or test tubes are commonly used in the chemical and biological arts to perform a variety of types of assays in a contained space. Assays that commonly have one or more steps performed in reaction receptacles include chemical reactions, immunoassays, and nucleic acid-based assays. Examples of such reactions and assays are thoroughly described in the available literature and are well known to those skilled in the art. While reaction receptacles are generally manufactured and sold as individual units or test tubes, it is common for practitioners to use holding racks to conveniently and collectively organize a group of reaction receptacles for performing multiple assays simultaneously or sequentially. In some instances, multiple reaction receptacles are assembled as a unitary piece.\nWith most assays, a substance transfer device is used to dispense solutions into or remove solutions from reaction receptacles. The most familiar substance transfer devices are pipettes and aspirators including one or more tubular elements through which fluids are dispensed or withdrawn. When substance transfer devices are used in conducting a group of independent assays at about the same time or in close proximity to one another, there is always the concern that a substance transfer device will inadvertently serve as a vehicle in transferring substances or contaminants between reaction receptacles. An additional concern is that the practitioner will improperly add substances into or remove substances from a reaction receptacle. To minimize the risk of cross-contamination and pipetting and aspirating errors, practitioners must carefully monitor substance transfers and exercise nearly flawless precision when pipetting substances into or aspirating substances from reaction receptacles. Avoiding cross-contamination and pipetting and aspirating mistakes is particularly important when the assay is diagnostic in nature or is designed to provide information concerning the progress of a patient's disease over time or the success of a treatment regimen.\nOne way to limit opportunities for cross-contamination is to reduce the amount of surface area on the substance transfer device that can come into contact with the contents of a reaction receptacle. This objective can be achieved by using a contact-limiting element, such as a pipette tip, which essentially serves as a barrier between the outer surface of the pipette and the contents of a reaction receptacle. And by selecting a pipette tip of sufficient length and volume, contact between the pipette and contents of a reaction receptacle can be substantially eliminated. This is because substances from the reaction receptacle will be drawn into a portion of the pipette tip which falls below the bottom surface of the pipette. Of course, in most instances, it will also be important to have a single pipette tip dedicated to each reaction receptacle.\nWhere a number of pipette tips are used to perform multiple assays simultaneously or sequentially, practitioners typically need to position a supply of pipette tips at a location that can be conveniently accessed by at least one pipette. Providing a sufficient quantity of pipette tips becomes more complicated when the substance transfer device functions robotically in an automated (or partially automated) assay instrument. In an automated format, a large reserve of pipette tips may need to be placed in the instrument at a site that is accessible by the pipette, but which limits the total amount of space required. Accordingly, there is a need for pipette tips that are readily accessible by a robotic pipette without requiring the pipette to engage in complicated movements or to travel over substantial distances.\nAnother problem presented by conventional reaction receptacles is that they come packaged as individual test tubes that are not amenable to manipulation by an automated assay instrument. Individual reaction receptacles hinder throughput efficiency since the practitioner and instrument must each handle the reaction receptacles separately. And because conventional reaction receptacles are not provided with any structure that permits them to be manipulated by an automated instrument, reaction receptacles are generally stationed at one situs within the instrument and are not afforded any automated mobility. This lack of movement imposes certain architectural limitations and assay inefficiencies since the instrument must be designed around the positioning of the reaction receptacles. Accordingly, there is a need for a reaction receptacle apparatus which can be manipulated by an automated assay instrument, where the apparatus may include one reaction receptacle or plurality of reaction receptacles coupled together as a single operative unit."}
{"text": "A problem associated with the burning of biomass fuel in general is the production of air pollutants. For example, the burning of biomass fuel (and particularly the inefficient burning of biomass fuel) may produce volatile, toxic, or other undesirable gases. Furthermore, large amounts of smoke and particulate matter may also be released into the atmosphere.\nIn this regard open fireplaces are particularly inefficient. That is, open fireplaces usually produce larger amounts of air pollutants, as compared to enclosed fireplaces. Furthermore, an open fireplace generally only provides heat directly in front of the fireplace, with the vast majority of the heat being lost up through the chimney or out through the rear wall of the fireplace.\nThe inefficiency of open fire places has been addressed to a certain extent by the use of domestic furnaces. Examples can be found in U.S. Pat. No. 4,559,882 (Dobson) and U.S. Pat. No. 4,630,553 (Goetzman).\nHowever, whilst the problems of the inefficient burning of biomass fuel for space heating can be addressed somewhat with a furnace, the extra capital cost is not always necessary, practical, or affordable. Furthermore, furnaces tend to be used mainly in very cold climates, but not in temperate climates, and are also usually coupled to some sort of central heating, which is not always desirable.\nMoreover, a further issue with furnaces is that they are usually completely closed off from view, and do not therefore provide the psychological or aesthetic benefits that are derived from lazy flame as a light source. That is, people like to see flames.\nPerhaps a result, “air tight wood burners” or simply “wood burners” have become increasingly more popular over the years, and are now in widespread use. “Wood Stove” is another common term for such appliances, particularly in North America.\nWood burners generally comprise a metal firebox, into which biomass fuel may be placed and burnt, an adjustable air control or damper, and an exhaust flue. Many, if not most, wood burners also have a glass door through which the fire and/or flames may be viewed.\nThe combustion of biomass fuel is a complex process, and includes a range of chemical reactions. As yet, there does not appear to be a wood burner or fireplace designed for space heating that adequately incorporates the various stages of combustion in relation to each other.\nThese stages are drying, pyrolysis, combustion and reduction, which if done correctly produce the combustible gases carbon monoxide and hydrogen. The carbon monoxide and hydrogen can then be combusted separately during what is known as secondary combustion to yield only water and carbon dioxide (and heat). However, most wood burners lack the ability to burn or convert such gases (and their precursors such as carbon dioxide and water) due to the wood burner not being able to produce enough heat to do so (conventional wood burners usually reach maximum temperatures of between 600° C.-800° C.).\nHence, attaining a high enough temperature during secondary combustion to consume all the volatiles distilled during pyrolysis or combustion is difficult because the necessary temperature is often higher than that which can ordinarily be generated. In this regard, the temperature required to adequately consume the vast majority (if not all) of the volatiles and/or particulate matter and smoke would be a minimum of approximately 900° C., and more preferably above 1000° C.\nU.S. Pat. No. 4,672,946 (Craver) describes a wood burner which has a secondary combustion means for burning the particulate matter in the flue gases. However, the temperature reached within the firebox of the Craver device is stated as being only around 540° C. (1000° F.) and the secondary combustion region only reaches up to around 760° C. (1400° F.). Hence, a disadvantage associated with Craver is that the design of the wood burner does not attain high enough temperatures to adequately consume the vast majority of gases or particulate matter. Furthermore, the wood burner described in Craver is not able to be retro-fitted to an existing wood burner or other fireplace.\nIn more recent times, many countries or local bodies have introduced regulations to restrict the sale of inefficient and/or polluting wood burners.\nFor example, in New Zealand the generally allowable standard for wood burners is a maximum of 1.0 grams of particulate matter released per kilogram of wood burned, accompanied by a minimum efficiency of 65%. However, some regions have gone further than this. For example, the Canterbury Regional Council in New Zealand (which is in the region of a weather-inversion layer) has lowered these levels to 0.6 grams of particulate matter per kilogram of wood burned. The Regulations further restrict the use of wood with a moisture content higher than 25%.\nHowever, these Regulations are not retrospective, and hence they only have effect in relation to wood burners manufactured and sold after the Regulations came into force. Moreover, to date there have been no innovations which have enabled people to bring their older wood burners up to modern compliance levels (voluntarily or otherwise).\nTwo factors which usually have the most detrimental effect regarding the efficiency of, and/or the release of air pollutants from, a wood burner are to do with refueling the wood burner and when shutting down or reducing the air supply to the wood burner.\nRefueling causes quenching, a situation where the introduction of fresh fuel to the fire is not supported by the heat contained within the existing fire to adequately pyrolyse the biomass. As a result, visible smoke and particulate matter are often seen exiting the top of the flue or chimney at this time. This can take a while to subside as enough heat builds up in the fire to commence the correct chemical processes required to efficiently combust the fresh fuel.\nA wood burner user may wish to reduce the air supply to keep the fire burning longer and/or while they are asleep. This is known as “banking”. In doing so, they generally place a full load of biomass fuel in the wood burner and shut down (or minimise) the air supply to prolong the burn time. However, the reduction in available oxygen and the corresponding detrimental effect on combustion results in more air pollutants being produced and released. Because this often results in the amount of air pollutants exceeding the minimum regulated amounts, many modern wood burner designs have denied the user the ability to shut down the air supply.\nThe air supply also affects the dynamics of wood burners because a greater draught causes more heat to be generated, but a greater portion of heat is lost up the flue. The higher velocity of gases also results in more particulate matter being exhausted to the atmosphere. Conversely, a lesser draught reduces the amount of particulate matter being drawn from the combustion chamber but also reduces the heat output. However it is possible in these conditions that although less heat is generated, less heat is also lost to the atmosphere as the heat has more time to radiate off before being exhausted.\nOr to put it another way, greater air means greater heat, but lower efficiency, however the greater heat actually results in a cleaner burn which lowers the emissions. With a lesser air supply, the fires get greater efficiency but the lower heat increases emissions. As a result of these dynamics, there is a common saying amongst laboratory engineers which is: “You can build a hot and clean fire, and you can build an efficient fire, but you can't build both in the same fire”. I believe that my combustion system, as described herein, does in fact result in both a hot and clean fire and an efficient fire.\nAnother approach taken by wood burner manufacturers to address problems of fire inefficiency or to reduce the release of pollutants is by employing the use of catalytic converters (either by retrofitting to existing wood burners or by incorporating them into new wood burners). However, catalytic converters are generally very expensive, and may be considered complex to operate and/or understand by many people—and this may be prohibitive to both wood burner manufacturers and end consumers. Furthermore, the installation of catalytic converters requires specialist knowledge and significant alterations to be made to the wood burner, and flue, and this can be time consuming, complex and expensive.\nIt may therefore be advantageous if there was available a relatively simple and/or improved combustion system, which included primary and secondary combustion zones which were able to result in more efficient combustion of biomass fuels and/or result in a lesser amount of air pollutants being released, as compared to presently available or prior art combustion systems.\nIt may also be of advantage if there was available a combustion system which was able to be retrofitted to existing fireplaces, such as wood burners, for example to increase their efficiency and/or to bring them up to modern compliance standards."}
{"text": "In five layers of cornea structure, endodermis located in the innermost layer has great significance on maintaining corneal transparency and normal corneal physiological function. Human Corneal Endothelial Cells (HCECs) regulate the transparency of cornea through pump function and barrier function. However, since adult corneal endothelial cells lack proliferation ability, its cytothesis mainly relies on cellular extension and migration of surrounding cells after damage. Corneal endothelial decompensation can be induced when the cell density is less than 500-800/mm2, leading to corneal edema and to losing transparency. Infection, inflammation, trauma, eye surgery and etc. may all induce endothelial cell damage, but due to the special physiological properties of human corneal endothelial cells, endothelial cell function recovery is often a difficult problem for eye disease treatment. The only effective way to treat corneal endothelial decompensation is to transplant corneal endothelial cells from a healthy donor and to replace the damaged or pathological changed cells by penetrating corneal transplantation or corneal endothelial transplantation. Given to worldwide extreme lack of available donor for corneal transplantation, and to postoperative immunological rejection, the application of corneal transplantation has been limited. Biological corneal transplantation can effectively relieve the pressure of natural corneal materials and has great value socially and economically. However, it demands sufficient supply of corneal endothelial cells for the purpose of the construction of biological corneal endothelial as well as merely endothelial cells transplantation. Due to limited proliferative capacity of adult corneal endothelial cells as well as the limited ability of external culture amplification, to find new endothelial cells is an urgent problem to be solved. It is extremely urgent to obtain a large number of functional coneal endothelial cells similar to normal human to meet the needs for cytobiology research and cell replacement therapy.\nAt present, research on the inducting differentiation of corneal endothelial cells is still in its early stage. Researchers earlier have tried to replace non-proliferable corneal endothelial cells with adipose skin cells or mesenchymal stem cells or tried to induce differentiation of embryonic stem cells, stromal stem cells, bone marrow endothelial progenitor cells or neural crest cells into corneal endothelial cells, but these methods all exist some problems, such as immune rejection, ethical issues, poor histocompatibility and lack of animal experiments, etc. A method for achieving a large number of functional corneal endothelial cells in vitro has not been found yet."}
{"text": "Vehicles are being provided with more complex systems. For example, vehicles commonly include a plethora of entertainment systems, such as stereos, USB interfaces for mobile telephones, video players, etc. Vehicles often have a host of other operator interfaces, such as emergency calling systems, vehicle navigation systems, heating and air conditioning systems, interior and exterior lighting controls, air bags, seatbelts, etc.\nVehicle operating environments are becoming more complex as well. For example, some roadways include u-turn lanes, round-a-bouts, no-left turn, multiple lanes one way in the morning and the other way in the afternoon, etc. Increases in traffic are also contributing to increased complexity.\nThese additional complexities contribute to increases in driver risk. What is needed are methods and systems for generating data representative of vehicle in-cabin insurance risk evaluations based on data representative of skeletal diagrams of a driver that are indicative of degrees of driver risk."}
{"text": "This invention relates to the detection of fluorescence and, in particular, the detection of fluorescence for cancer diagnosis.\nWith conventional endoscopic techniques, early cancer detection presents a diagnostic challenge even for experienced physicians. In its early stages, a cancerous tumor may only be a few cell layers thick, and thus, the surface of the tissue may not physically change sufficiently to be visible when viewed under \"white light\" illumination with conventional endoscopic techniques. By \"white light\" illumination, is meant a relatively broad range of the visible portion of the electromagnetic spectrum.\nFluorescence imaging techniques have been used to detect the presence of cancer. Fluorescence is generally defined as the emission of radiation, usually in the form of visible light, resulting from the absorption of radiation in another spectral band from a radiation source. The impetus for using fluorescence imaging in cancer detection was due in part to the recognition that normal and malignant tissue provide different fluorescence features.\nOne approach for using fluorescence imaging in cancer detection involves comparing the intrinsic fluorescence (also referred to as \"autofluorescence\" or background fluorescence that is generated by all cells due to the natural presence of fluorescent molecules inside them) of normal tissue with the intrinsic fluorescence of malignantly-transformed tissue.\nAnother approach for fluorescence imaging in cancer detection involves the use of exogenic substances, including fluorescent dyes, which are administered into a patient and accumulate in malignant tissue. Experiments with cultured cells indicate that the primary cause of dye accumulation in malignant tissue is due to the different metabolism rate of normal and cancerous cells. One advantage of this approach is that the fluorescence can be detected at a peak of the dye's fluorescence spectrum."}
{"text": "It is known to have an engine having an engine case, wherein on an upper surface of the engine case (e.g., which accommodates a crankshaft, a clutch mechanism, and a transmission mechanism) there is provided an oil passage controlling apparatus that controls switching of oil passages so as to control the operations of other components (e.g., the transmission mechanism). Such an engine structure is known to include a cover which covers the oil passage controlling apparatus (for example, refer to JP-A No. S63-125490.)\nOne problem with this known structure is that the provision of a special cover for covering the oil passage controlling apparatus (the apparatus that controls the operations of the other components) increases both weight and cost (e.g., due to an increase in the number of components.) Another problem is that the oil passage controlling apparatus that controls the operations of the other components using the oil passages is located outside the engine case; thus hydraulic pathways provided between the oil passage controlling apparatus and the engine case cross over the engine case. Accordingly, high sealing performance is required between the oil passage controlling apparatus and the engine case, which has an impact on the attachability of the oil passage controlling apparatus.\nWith respect to the apparatus that controls the operations of the other components, particularly in an engine including clutch actuators which control a connecting/disconnecting operation of a clutch mechanism, one object of the present invention is to reduce the number of components and to improve the attachability of the clutch actuators by eliminating a cover for the clutch actuators."}
{"text": "The majority of deaths associated with malignant tumor are due to the metastasis of the original tumor cells to tissues and organs distant from the initial tumor. Accordingly, early diagnosis of metastasis is a critical factor for the survival of a cancer patient, and early diagnosis of tumor and monitoring of tumor growth are considered as very important factors for successful treatment of a cancer patient. The diagnosis of cancer usually uses diagnosis techniques by histopathology. The histopathological diagnosis technique is a method of using a tissue sample from a living subject to diagnose a cancer. Such a histopathological approach allows a tumor cell to be directly observed. However, it may be incorrect whether there is a tumor from a tissue site selected in order to obtain a sample from a living subject, and only data about a particular site obtained from the living subject are provided and thus it is difficult to know whether tumor has metastasized to another site. For this reason, the applicability in diagnosing and monitoring tumors may be limited.\nIt is known that circulating tumor cells (CTCs) are found from a patient before the tumor is originally detected. Accordingly, CTCs may play an important role in early diagnosis and prognosis of cancers. In addition, because cancer usually metastasizes through blood, CTC may be a marker for determining whether cancer has metastasized. Even after cancer cells have been removed by surgery, CTCs may be still exist and cancer may reoccur. However, very small amounts of these CTCs are found in blood and the cells are themselves weak, and thus it is very difficult to detect them and grasp the number of the cells. Accordingly, there still remains a need for a diagnosis method that is highly sensitive to detect CTCs, cancer cells, or cancer stem cells in a patient's body.\nCTC separation methods by using magnetic nanoparticles are described in the related art. However, the method according to the related art is disadvantageous because the processes are very complicated, for example, the method of separating serum from blood, and using the affinity of biotin and streptavidin even in a magnetic separation process have a risk of losing of CTCs in the separation steps.\nAccordingly, there still remains a need for a method for efficiently separating tumor cells from a biological sample and an apparatus associated with that."}
{"text": "The present invention relates to an assembly for cooling at least one electronic component disposed on a printed circuit board in accordance with the type defined in greater detail in the preamble of claim 1.\nBy way of example, a device for cooling an electronic device and an associated production method are known from DE 10 2013 204 029 A1. The cooling device is intended for a control device, which is disposed in a housing and comprises an electronic circuit disposed on a printed circuit board, having an element that produces a lossy line, wherein a contact pressure for establishing mechanical and thermal contact with the printed circuit board is provided on the housing. The contact pressure is disposed adjacent to the element producing a lossy line.\nTensions arise due to the different heights and tolerances of the electrical components in the circuitry, and because of bends in the printed circuit board, or carrier plate, respectively, which lead to the cooling surfaces of the components or printed circuit board not lying, or lying with different pressures, against the cooling device. As a result of the high working temperatures, the poorly cooled elements reduce the lifetime of the electronic components."}
{"text": "This invention relates generally to an apparatus for the machining of a tubular workpiece, and more particularly to such an apparatus having a work implement, such as a metal cutting burner or a metal welding tool, capable of a wide variety of unobstructed movements relative to the workpiece.\nThe free end of the workpiece implement is capable of being longitudinally adjusteed along its central axis, the work implement is mounted on a support arm for movement therealong, the support arm is mounted on the apparatus for movement about a vertical swivel axis lying perpendicular to the support arm, and the work implement is capable of movement about a tilting axis lying at a right angle to the swivel axis.\nA known apparatus of the general type aforedescribed is set forth in a publication entitled: \"RB-Elektron 650/1200 AGD\". The support arm of such apparatus is cantilevered horizontally from a support frame capable of being moved in a vertical direction and in a direction of the longitudinal axis of the tubular work-piece. The support arm is mounted on the frame for movement about a swivel axis which lies horizontally. Tilting movement of the burner is effected by parallelogram guiding. By swivelling the burner arm about its horizontal swivel axis and tilting the burner about the tilting axis of the burner arm, the burner is capable of being adjusted throughout almost 360.degree.. Only angular areas of adjustment with approximately a horizontal burner position are excluded since such positions cannot be achieved, due to the geometry of the parallel guiding. In other words, from a central vertical position, the burner cannot be swivelled completely through .+-.90.degree. therefrom. And, the central point of the angular adjustability of the burner is located by almost an angle of 360.degree. on the intersection of the swivel and burner axes.\nThis known burner apparatus presents drawbacks in that the burner arm has too much play whih can only be corrected by additional and extensive mechanical arrangements. The parallelogram guiding system, necessarily includes a large number of joints which must be capable of movement without interference so that the tilting may be carried out with essentially the same ease independent of the initial tilting angle, and which must be adjusted in such a precise manner that, when the burner is swivelled through its vertical position, no dead angles occur. These requirements cannot be achieved in practice with the known burner arm arrangement. Instead, burner arm swivelling becomes unstable during repeated and extensive use and, in practice, it makes jerk-like movements since the tolerances of the swivel angles are too high. Thus, precise cutting requirements are not capable of being achieved."}
{"text": "The invention relates to boltless pipe connectors suitable for attachment on or adjacent the end of a plain end pipe to enable coupling of said pipe to a further pipe, pump, valve, end cap or the like. These connectors are universal, that is they have a wide tolerance within a certain size range: pipes of a given nominal diameter do not all have the same outside diameter.\nEP 727026 uses grease under pressure to grip and seal pipes as they are joined. This technology has been limited to small tolerances, typically 4 or 5 mm, in pipe diameter. The maximum movement of the grippers is approximately half thickness of the grippers or they would slip out of the container groove in the housing.\nIn the connector described and shown in the present application, a seal located in an annular recess in an inner surface of the connector has segmental grippers of L-shaped cross-section comprising a leg between the seal and a wall of the recess, and a foot having a serrated inner surface for engagement with the pipe. This allows for a tolerance exceeding 10 or 15 mm thereby allowing the pipe connection to be used with all pipes of a given nominal diameter.\nA metal shim may be provided between the legs of the adjacent grippers. Grippers may be provided between one or each side of the seal and an adjacent wall of the recess. The housing is provided with an inlet for a sealing fluid which may be a nipple for the supply of grease or oil under pressure or components for the formation of an epoxy resin to make the seal permanent. The housing is generally provided with an air bleed offset from the sealing fluid inlet for use during assembly when the recess is filled with grease or oil as sealing fluid."}
{"text": "Biological preparations, for example in the form of plant specimens must be prepared before they can be inspected under a microscope, such as an electron microscope. The preparation requires subjecting the specimens to a fixation process so that the tissue will retain its structure without deterioration prior to and during the microscopic inspection. The fixation process uses a fixation substance such as an aldehyde that penetrates into the tissue of the specimen, thereby preventing a change in the specimen that could otherwise be caused by the subsequent embedding steps and by the vacuum prevailing in the specimen chamber of an electron microscope.\nConventional fixation substances include, among others, aldehyde, primarily glutaraldehyde, formaldehyde, or even the unsaturated aldehyde acrolein, whereby the fixation substances are frequently also used as mixtures of these formaldehydes.\nConventionally, fixation substances are used in the form of aqueous solutions for the so-called immersion fixation. Less frequently, fixation may also be accomplished by using these substances in the gas phase for performing a so-called fixation by fumigation.\nA substantial problem has been encountered heretofore in conventional fixation processes with regard to the speed at which the fixation substance enters into the tissue and fixes the fine structure of the specimen prior to any possibility of a change in that fine tissue structure. Such changes must be avoided since they defeat the purpose of the microscopic inspection. Especially the preparation of plant specimens or rather their fixation poses two problems. On the one hand, above ground plant portions are enclosed by a substantially water impermeable outer skin, the so-called cuticula. On the other hand, air is present in the intercellular interstices in the leaves of plants, whereby the penetration of liquid, including fixation solutions into the leaf tissue is made more difficult.\nIn a book entitled \"Basics of Histochemistry\" (Grundlagen der Histochemie), by Hans Luppa, published by Vieweg Verlag Braunschweig, 1978, section 1, pages 42 and 43, a method has been described for fixing biologic preparations by using aldehyde fixation substances. This publication further describes on page 52 that trichloroacetic acid may be effective by hydrolysis on the heteropolar valence bindings to thereby loosen the hydration status and that trichloroacetic acid penetrates quickly into the tissue to thereby cause a type of preliminary fixation.\nIn connection with most of the conventionally used fixation methods, the plant specimens are first mechanically comminuted. Thereafter, different fixation sequential steps are employed, depending on the particular type of plant specimen. These conventional fixation steps require a relatively large volume throughput of fixation solution, the disposal of which subsequently causes a critical environmental problem. Further, conventional fixation methods require repeated manual intervention by the operator so that conventional methods are not suitable for an automatically performed fixation sequence. Such automatically performable sequence of fixation steps, however, is indispensable for experiments to be performed, particularly in a spacecraft in outer space or in areas that are difficult to access, for example, outside in the field."}
{"text": "In the past decade, the science of genetic engineering has developed rapidly. A variety of processes are known for inserting a heterologous gene into bacteria, whereby the bacteria become capable of efficient expression of the inserted genes. Such processes normally involve the use of plasmids which may be cleaved at one or more selected cleavage sites by restriction endonucleases. Typically, a gene of interest is obtained by cleaving one piece of DNA, and the resulting DNA fragment is mixed with a fragment obtained by cleaving a vector such as a plasmid. The different strands of DNA are then connected (“ligated”) to each other to form a reconstituted plasmid. See, for example, U.S. Pat. Nos. 4,237,224 (Cohen and Boyer, 1980); 4,264,731 (Shine, 1981); 4,273,875 (Manis, 1981); 4,322,499 (Baxter et al, 1982), and 4,336,336 (Silhavy et al, 1982).\nA variety of other reference works are available. Some of these works describe the natural process whereby DNA is transcribed into mRNA and mRNA is translated into protein, see, e.g., L. Stryer, Biochemistry, 2nd edition, p 559 et seq. (W. H. Freeman and Co., 1981); A. L. Lehninger, Biochemistry, 2nd edition, p. 853 et seq. (Worth Publ., 1975). Other works describe methods and products of genetic manipulation; see, e.g., T. Maniatis et al, Molecular Cloning, A Laboratory Manual (Cold Spring Harbor Labs, 1982); J. K. Setlow and A. Hollaender, Genetic Engineering, Principles and Methods (Plenum Press, 1979). Hereafter, all references will be cited in abbreviated form; a list of complete citations is included after the Examples.\nMost of the genetic engineering work performed to date involves the insertion of genes into various types of cells, primarily bacteria such as E. coli, various other microorganisms such as yeast, and mammalian cells. However, many of the techniques and substances used for genetic engineering of animal cells and microorganisms are not directly applicable to genetic engineering involving plants.\nAs used herein, the term “plant” refers to a multicellular differentiated organism that is capable of photosynthesis, such as angiosperms and multicellular algae. This does not include microorganisms, such as bacteria, yeast, and fungi. The term “plant cell” includes any cell derived from a plant; this includes undifferentiated tissue such as callus or crown gall tumor, as well as plant seeds, propagules, pollen, or plant embryos.\nTi and Ri Plasmids\nThe tumor-inducing (Ti) plasmid of Agrobacterium tumefaciens has been proposed for use as a natural vector for introducing foreign genetic information into plant cells (Hernalsteen et al 1980; Rorsch and Schilperoort, 1978). Certain types of A. tumefaciens are capable of infecting a wide variety of plant cells, causing crown gall disease. The process of infection is not fully understood. At least part of the Ti plasmid enters the plant cell. Various metabolic alterations occur, and part of the Ti plasmid is inserted into the genome of the plant (presumably into the chromosomes). The part of the Ti plasmid that enters the plant genome is designated as “transferred DNA” (T-DNA). T-DNA is stably maintained in the plant DNA (Chilton et al, 1977; Yadev et al, 1980; Willmitzer et al, 1980; Otten et al, 1981).\nResearch by several laboratories has led to the characterization of several structural (i.e., protein coding) genes located in T-DNA (Garfinkel et al 1981; Leemans et al 1982), as well as other DNA sequences which appear to serve various other functions. For example, two sequences referred to as the “left border” and the “right border” appear to delineate T-DNA and may be involved in the process whereby T-DNA is transferred into plant chromosomes (Zambryski et al 1982).\nA different species of Agrobacterium, A. rhizogenes, carries a “root-inducing” (Ri) plasmid which is similar to the Ti plasmid. Infection of a plant cell by A. rhizogenes causes hairy root disease. Like the Ti plasmid, a segment of DNA called “T-DNA” (also referred to by some researchers as “R-DNA”) is transferred into the plant genome of an infected cell.\nVarious other bacteria are also reported to be capable of causing genetic transformation of plant cells, including A. rubi and certain bacteria of the genus Rhizobia which have been treated with a mutagenic agent. Hooykaas et al, at page 156 of Setlow and Holaender, 1979.\nAs used herein, the term “plant tumor inducing plasmid” includes any plasmid (1) which is contained in a microorganism, other than a virus, which is capable of causing genetic transformation of one or more types of plants or plant cells, and (2) which contains a segment of DNA which is inserted into a plant genome. This includes Ri plasmids.\nAs used herein, the term “T-DNA” to a segment of DNA in or from a plant tumor inducing plasmid (1) which has been inserted into the genome of one or more types of plant cells, or (2) which is contained in a segment of DNA that is located between two sequences of bases which are capable of serving as T-DNA borders. As used herein, the terms “T-DNA border” and “border” are determined and applied empirically; these terms shall refer to a sequence of bases which appears at or near the end of a segment of DNA which is transferred from a plant tumor inducing plasmid into a plant genome.\nDespite the existing knowledge of T-DNA and plant tumor inducing plasmids, no one prior to this invention has been able to utilize these vectors for the introduction of foreign genes which are expressed in genetically modified plants. A variety of obstacles to such use have been encountered in genetic engineering efforts. Such obstacles include:\n1) the large size (approximately 200,000 base pairs) and resulting complexity of Ti plasmids preclude the use of standard recombinant DNA techniques to genetically modify and/or insert foreign genes into specific sites in the T-DNA. For example, there are no known unique restriction endonuclease cleavage sites in a plant tumor inducing plasmid (Leemans et al, 1982).\n2) the T-DNA, which is inserted into and expressed in plant cells, contains genes which are involved in the production of high levels of phytohormones in the transformed plant cells (Leemans et al 1982). The high levels of phytohormones interfere with the normal metabolic and regenerative process of the cells, and prevent the formation of phenotypically normal plants from the cells (Braun and Wood, 1976; Yang et al, 1980). Exceptions to this are rare cases where the T-DNA has undergone extensive spontaneous deletions in planta to eliminate those genes involved in phytohormone production. Under these conditions, normal plants are reported to be obtainable at low frequency (Otten et al, 1981). However, the T-DNA genes involved in phytohormone production could not be deleted prior to this invention, since they were very important in the identification and/or selection of transformed plant cells (Marton et al, 1979).\nAs described above, simple recombinant DNA techniques for introducing foreign genes into plasmids are not applicable to the large Ti plasmid. As a result, several indirect methods have been developed and are discussed below. The first reported use of the Ti plasmid as a vector was in model experiments in which bacterial transposons were inserted into T-DNA and subsequently introduced into plant cells. The bacterial transposons were reported to be stably maintained in the plant genome (Hernalsteens et al, 1980; Garfinkel et al 1981). However, in these cases the transformed tumor tissues were found to be incapable of regeneration into normal plants, and there was no reported evidence for the expression of bacterial genes in the plant cells. In addition, because the insertion of bacterial transposons is believed to be essentially random, a great deal of effort was required to identify and localize the position of the inserted DNA in these examples. Therefore, this approach is not likely to be useful to introduce desired genes in a predictable manner into plants.\nOther researchers have reported the use of intermediate vectors which replicate in both E. coli and A. tumefaciens (Matzke and Chilton, 1981; Leemans et al 1981; Garfinkel et al, 1981). The intermediate vectors contain relatively small subfragments of the Ti plasmid which can be manipulated using standard recombinant DNA techniques. The subfragments can be modified by the deletion of specific sequences or by the insertion of foreign genes at specific sites. The intermediate vectors containing the modified T DNA subfragment are then introduced into A. tumefaciens by transformation or conjugation. Double recombination between the modified T-DNA fragment on the intermediate vector and its wild-type counterpart on the Ti plasmid results in the replacement of the wild-type copy with the modified fragment. Cells which contain the recombined Ti plasmids can be selected using appropriate antibiotics.\nVarious foreign DNA's have been inserted at specific sites in the T-DNA by this method and they have been reported to be stably transferred into plant cells (Matzke and Chilton, 1981, Leemans et al 1981, 1982). However, such foreign genes have not been reported to be capable of expression in plant cells, and the transformed plant cells remain incapable of regeneration into normal plants. Furthermore, in the procedure described above, it is preferred for a double crossover event to occur, in order to substitute the modified DNA fragment for the wild-type copy. A single crossover results in the formation of a co-integrate plasmid which contains two copies of the T-DNA subfragments. This duplication is undesirable in these methods since homologous recombination, which can occur in A. tumefaciens cells or in plant cells, can result in the loss of the inserted foreign gene(s).\nA major disadvantage of the above methods is that the frequency of double recombination is quite low, about 10−4 to 10−9 (Leemans et al, 1981) and it requires extensive effort to identify and isolate the rare double-crossover recombinants. As a result, the number and types of experiments which can be performed using existing methods for genetically engineering the Ti plasmid is severely limited.\nOther Means for Inserting DNA into Plant Cells\nA variety of other methods have been reported for inserting DNA into plant cells. One such method involves the use of lipid vesicles, also called liposomes, to encapsulate one or more DNA molecules. The liposomes and their DNA contents may be taken up by plant cells; see, e.g., Lurquin, 1981. If the inserted DNA can be incorporated into the plant genome, replicated, and inherited, the plant cells will be transformed.\nOther alternate techniques involve contacting plant cells with DNA which is complexed with either (a) polycationic substances, such as poly-L-ornithine (Davey et al, 1980), or (b) calcium phosphate (Krens et al, 1982). Using these techniques, all or part of a Ti plasmid has been reportedly inserted into plant cells, causing tumorigenic alteration of the plant cells.\nAnother method has been developed involving the fusion of bacteria, which contain desired plasmids, with plant cells. Such methods involve converting the bacteria into spheroplasts and converting the plant cells into protoplasts. Both of these methods remove the cell wall barrier from the bacterial and plant cells, using enzymic digestion. The two cell types can then be used together by exposure to chemical agents, such as polyethylene glycol. See Hasezawa et al, 1981.\nHowever, all of the foregoing techniques suffer from one or more of the following problems:\n1. transformation efficiencies reported to date have been very low;\n2. only small DNA molecules can be inserted into plant cells;\n3. only small numbers of DNA molecules can be inserted into plant cells; and/or,\n4. a gene which is inserted into a plant cell will not be stably maintained by the plant cell unless it is incorporated into the genome of the plant cell, i.e., unless the gene is inserted into a chromosome or plasmid that replicates in the plant cell.\nFor these and possibly other reasons, no one has yet reported expression of a gene inserted into a plant cell by any of the foregoing techniques, except for the tumorigenic transformations noted above.\nPrior to this invention, no satisfactory method existed for the creation and identification of genetically transformed plant cells which could be routinely regenerated into morphologically normal plants."}
{"text": "Modern pick-up trucks and similar vehicles commonly have a second or subsequent row of one or more seats positioned rearwardly of the driver's seat. Such second or subsequent rows of seats are commonly movable so as to provide an increased cargo area in the vehicle's passenger compartment or “cabin”. Such second and subsequent rows of seats are hereinafter referred to alternately as “second row seats” and/or as “second row seat assemblies”.\nSecond and subsequent row seats for pick-up trucks and similar vehicles are typically of the general type having a seat portion and a backrest portion. In the prior art, such second and subsequent row seats have been provided with an in-use or deployed design configuration, whereat the seat portion thereof is substantially horizontal and either touches or is in very close relationship with the rear wall of the vehicle's cabin, and have been rearward pivotable therefrom towards a raised stowed configuration, whereat the seat portion is substantially vertical and substantially adjacent to the cabin rear wall. One shortcoming with such second row seat designs according to the prior art, however, is that thus far, in the deployed design configuration, they have been insufficiently spaced from the rear wall of the cabin to allow for sufficient angular adjustment of their respective backrest portions. As such, second row seat passengers in pick-up trucks and the like have heretofore been faced with the prospect of riding in significant discomfort in a seat with a backrest oriented substantially upright.\nManufacturers of mini-vans and sport utility vehicles have previously provided second and subsequent rows of seats that are spaced from the rear wall of the passenger compartment, whilst also being pivotally movable between a deployed design configuration and a raised stowed configuration. U.S. Pat. No. 4,700,989 (Ercilla) discloses a vehicle seat assembly that is forwardly pivotable from a deployed design configuration towards a raised stowed configuration. In the deployed design configuration, the legs of the Ercilla vehicle seat assembly are locked to fixed attachment rods supported by the vehicle floor. The front legs of the seat disclosed by Ercilla include a blocking mechanism that prevents unlocking of the front locking means when the vehicle seat assembly is in its deployed design configuration. In order to pivot the vehicle seat assembly from its deployed design configuration to its raised stowed configuration, a locking means at the rear legs is released, and then the vehicle seat assembly is pivoted forwardly to the raised stowed configuration. Ordinarily, the front locking means remains engaged, but it may be separately released when the seat is in the raised stowed configuration, if the Ercilla vehicle seat assembly is to be entirely removed from the vehicle.\nA common problem associated with all seats of the general type having a raised stowed configuration that is located forwardly of the deployed design configuration, however, is that any increased cargo space provided in the passenger compartment of the vehicle when such a seat assembly is in its raised stowed configuration is not easily accessible other than by way of a rearward-facing door or hatchback. That is, users attempting to load or unload cargo from such vehicles through side-facing doors will usually encounter significant difficulties, especially insofar as any second row seats positioned in the raised stowed configuration may present a serious obstacle to such effort by partial occlusion of said side-facing doors. It will, of course, be appreciated that most pick-up trucks are not provided with rearward-facing doors or hatchbacks. As such, and because any cargo to be carried in the cabs of pick-up trucks must therefore be loaded or unloaded through the vehicle's side-facing doors, forwardly pivoting stowable second row seats have not been particularly useful in this context.\nOver and above any of the aforesaid difficulties that may be encountered in adapting seat assemblies disclosed by Ercilla for a particular use as second row seats in the cabs of pick-up trucks and the like, it is also to be noted that the Ercilla patent fails to disclose any means for supporting the vehicle seat assembly in its raised stowed configuration. However, while the Ercilla patent does not itself disclose a means for preventing unexpected movement of the seat assembly towards the deployed design configuration during acceleration of the vehicle so as to thereby avoid damage to any cargo stored therewithin, the production version of this device actually uses a strap connected between the seat portion and the mini-van roof for this purpose. The strap provided in the production version of this device, however, is quite awkward to employ and aesthetically unpleasing, whilst also potentially posing a serious safety problem in crash-like situations, where the loading on such a strap can be unusually high. Moreover, it should also be appreciated that strap arrangements permit an inherent degree of pivotal movement of the raised stowed seat in at least one direction, which movement results in unacceptable movement, squeaks, rattles and other noise during travel of the vehicle over rough terrain or other surfaces. Thus, among other things, it will be seen to be desirable to rigidly secure a second row vehicle seat assembly in the raised stowed configuration in a manner which prevents any significant amount of motion of the second row vehicle seat assembly in this position. Moreover, it is desirable that a second row vehicle seat assembly be rigidly securable in a manner which is easy and convenient to use.\nIt is, therefore, one object of the present invention to provide a second row vehicle seat assembly that is pivotally movable between a deployed design configuration and a raised stowed configuration.\nIt is a further object of the present invention to provide a second row vehicle seat assembly which, in the raised stowed configuration, is positioned substantially adjacent to the rear wall of the vehicle's cabin.\nIt is another object of the present invention to provide a second row vehicle seat assembly that provides a stowed configuration which maximizes cargo space forwardly of the vehicle seat assembly.\nIt is yet another object of the present invention to provide a second row vehicle seat assembly that, in its deployed design configuration, is sufficiently spaced from the rear wall of the cabin so as to provide significant additional cargo volume therebetween and so as to provide ample space therebetween to adjust the inclination angle of the seat's backrest portion to comfortably accommodate a seat occupant or user.\nIt is a still further object of the present invention to provide a second row vehicle seat assembly having a seat cushion member that is selectively releasable from secure engagement with an interior floor of the cabin only when the seat assembly is in an unoccupiable fold-flat configuration.\nIt is a yet still further object of the present invention to provide a second row vehicle seat assembly that has a floor release mechanism that is convenient to reach and use in the unoccupiable fold-flat configuration.\nIt is still another object of the present invention to provide a second row vehicle seat assembly that may be locked in the raised stowed configuration, such that the raised stowed seat assembly is rigidly secured against pivotal movement upon forward and rearward acceleration or deceleration of the vehicle in which it is installed.\nIt is a further object of the present invention to provide a second row vehicle seat assembly that may be easily and conveniently locked in the raised stowed configuration."}
{"text": "The present invention relates to model aircraft and in particular to the remote radio control of such aircraft by means of series of PDM signals.\nOne of the major causes of accidents involving model airplanes is loss of control. Each year deadly accidents occur as a result of persons being struck by such models or slashed by the rotating propeller blades of a model out of control.\nTests conducted in Germany show that in 80% of model airplane accidents the aircraft remote control apparatus was intact after the accident indicating that the malfunction that lead to the accident was caused by interference with control signal transmission. The loss of communication between the model operator on the ground and his aircraft is thus believed to be a major cause of model aircraft accidents.\nIn view of the above, it is the principal object of the present invention to provide a system for the safe control of model aircraft in the event of a loss of communication between a model operator and his aircraft.\nA further object is to provide such a system which automatically shifts the model into a predetermined pattern in the event of a loss of communication which permits the model to be safely brought to the ground."}
{"text": "The following includes information that may be useful in understanding the present invention(s). It is not an admission that any of the information provided herein is prior art, or material, to the presently described or claimed inventions, or that any publication or document that is specifically or implicitly referenced is prior art.\n1. Field of the Invention\nThe present invention relates generally to the field of support structures for safely carrying weapons and more specifically relates to a gun mounting apparatus for utility terrain vehicles.\n2. Description of the Related Art\nThe use of ATVs and UTVs has become increasingly popular. Given that weapons are often used in outdoor activities, many people install gun racks or gun mounts to their ATVs or UTVs to safely transport their weapons. Weapon racks and mounts for vehicles' prime function is to safely secure and transport weapons. The use of weapon racks and mounts for vehicles may further protect weapons against damage and accidental discharge which may occur if a weapon is not secured during transportation. Weapon racks and mounts should also facilitate uncomplicated, rapid, and hassle free access to the weapon.\nVarious gun racks and mounts have been proposed for vehicles. Unfortunately, many traditional gun racks that are currently available may be cumbersome and unsafe. Many may take up space in the personal cabin of vehicles, causing passengers to shift uncomfortably to avoid the guns. This also threatens the safety of passengers, as the guns or racks may strike them while vehicles are in motion. Many gun racks may be difficult to install, and may not allow for efficient application and removal of guns. Many gun racks and mounts comprise many parts, making it complicated to install or uninstall them from vehicles and to store them.\nTherefore a need exists for a gun mount for a vehicle which comprises few parts, is easy to install, keeps the weapon in a safe and comfortable position for driver and passenger and which simultaneously provides for easy access to the weapon and protects the weapon from accidental discharge and from being scratched or damaged.\nVarious attempts have been made to solve the above-mentioned problems such as those found in U.S. Pat. Nos. 5,495,969, 3,876,078, 2,251,271, 6,986,446, 7,137,511, and 4,995,537. This prior art is representative of support structures for safely carrying weapons. None of the above inventions and patents, taken either singly or in combination, is seen to describe the invention as claimed.\nIdeally, a gun mounting apparatus for utility terrain vehicles should be easy to assemble, comprise few parts, be easy to install and, yet, operates reliably and be manufactured at a modest expense. Thus, a need exists for a reliable UTV Gun Mount System to securely store and transport weapons and to avoid the above-mentioned problems."}
{"text": "1. Field of the Invention\nThe present invention relates to a projection display apparatus for projecting an image on a projection plane.\n2. Description of the Related Art\nRecently, the development of projectors in which lasers with large radiation energy are used as light sources has been under way. It is necessary to take sufficient measures for a person not to erroneously enter a projection space through which the light that has been projected from such a projector pass.\nAccordingly, a technique of detecting an entering person by a sensor, etc., has been proposed.\nThe aforementioned technique of detecting an entering person by using a sensor is effective when a screen is installed at a normal position and in a normal state. However, if a screen is not installed at a normal position or if a hole or a hollow is present in a screen, there is a possibility that the projected light may be guided to an unintended area."}
{"text": "As known in the art, wireless communications networks are typically formed from a plurality of basestations (also referred to as base transceiver stations, BTSs), each basestation being located in a cell which is a geographical region served by that basestation. User terminals, also referred to as user equipments (UEs), are able to communicate with one or more of the basestations using wireless communications.\nPreviously many wireless communications networks have used so called “hard handoff” schemes whereby an uplink signal from a UE to a basestation is only actioned at one of the basestations. That is, the uplink signal is only demodulated and decoded at a single “best” cellsite or basestation. Thus, the uplink signal from a UE only results in communication at one of the basestations. As a UE moves, for example, a mobile phone in a car that is travelling, the UE may move between cells and be “handed off” from uplink communication with one basestation to uplink communication with another basestation.\nIt has also been recognised that uplink signals from one UE are often received by a plurality of different basestations (or sectors at the same basestation) in a wireless communications network, and that by making use of the signals received at two or more of those basestations (or sectors), improved uplink communication is achieved. For example, FIG. 1 is a schematic diagram of a wireless communications network comprising a plurality of basestations 10 each within a tri-sectored cell 11. (In FIG. 1 the sectors are represented by hexagonal shapes.) A plurality of UEs are shown, each represented by a dot 12. Two uplink signal paths from most of these UEs 12 are represented by lines, each to a different basestation 10. The darker of the two lines is intended to represent the uplink signal path giving the stronger signal received at the basestation.\nThus more recently, “soft handoff” schemes have been used and are part of the IS95 standard. In such “soft handoff” schemes the uplink signal is separately demodulated and decoded on a frame-by-frame basis at multiple basestations. The hard-decoded frames (i.e. containing a 1 or 0 in each bit position) from these basestations are then backhauled (communicated in some way) to a common point. This is illustrated schematically in FIG. 2a which shows three basestations 20 each in a tri-sectored cell. A user equipment, in this case a mobile telephone handset 21 is shown with four uplink communication paths indicated by arrows 1 to 4. Arrows 1, 2 and 4 show paths to different basestations whereas arrows 1 and 3 show paths to different sectors of the same basestation. The degree of soft handoff is said to be 4-way because four different uplink paths are used. Each of the basestations demodulates and decodes the uplink signal received from the UE 21. The resulting hard-decoded frames are then backhauled to a common point 22 as indicated by arrows 23, 24, 25.\nFor example, the common point is typically a selector bank subsystem (SBS) of a base station controller (BSC). A processor at the common point then selects the “best” of the hard-decoded frames received from the different cell sites and discards the rest. For example, the SBS chooses as the single survivor frame (for vocoder decode) one of the backhauled set (if possible) with a correct cyclic redundancy check (CRC) (i.e. a CRC which correctly corresponds to the decoded data). In this way improved uplink communication is achieved (as compared with hard handoff) but at the expense of carrying out extra processing at the basestations, providing and carrying out the backhaul communication and carrying out the additional processing at the common point.\nmDiversity Inc. describe a variation on this soft handoff scheme in their U.S. Pat. Nos. 6,349,094 and 6,192,038. They describe a plurality of collectors (similar to basestations) distributed at macro-diverse locations for receiving reverse channel (uplink) signals from users. Each of the collectors processes the uplink signals to yield one or more sequences of data bits as a burst and corresponding initial confidence metrics for each bit. The collectors forward the data bits and confidence metrics to an aggregator or common point. At the aggregator the data for the same user received from the multiple collectors is combined by using the confidence metrics to make a final decision on each bit.\nOne problem with previous types of soft-handoff methods is that relatively complex processing is required at each basestation in order to demodulate and decode the signals.\nIt is also desired to further increase uplink capacity as compared with known types of soft-handoff methods.\nAn object of the present invention is to provide an improved soft-handoff method which overcomes or at least mitigates one or more of the problems noted above.\nFurther benefits and advantages of the invention will become apparent from a consideration of the following detailed description given with reference to the accompanying drawings, which specify and show preferred embodiments of the invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device and a method of manufacturing the same and, more particularly, a semiconductor device having a plurality of wiring formed to space apart from each other and a method of manufacturing the same.\n2. Description of the Prior Art\nIn semiconductor integrated circuit devices, multilayered wiring and high density wiring have been required according to high integration of the device.\nSince an uneven surface of an interlayer insulating film is readily formed in the multilayered wiring structure, a wiring formed on the interlayer insulating film must be prevented from being cut off by reducing differences in the unevenness. In general, as a method of preventing the cut-off of the wiring, the upper surface of the interlayer insulating film has been planarized.\nUsually planarization of the interlayer insulating film is effected by following steps.\nAs shown in a plan view of FIG. 1A, plural lower wiring 3 are first formed in parallel on a base insulating film 2 formed on a surface of a semiconductor substrate 1. Then, an SiO.sub.2 film 4 serving as the interlayer insulating film is formed on an entire surface of the semiconductor substrate 1 by CVD process. In this case, as shown in a sectional view of FIG. 1B, the SiO.sub.2 film 4 is grown on both surfaces of the base insulating film 2 and the wiring 3. For this reason, an uneven surface shape is generated on the SiO.sub.2 film 4 according to the surface unevenness of the base insulating film 2 and the wiring 3.\nSubsequently, silicon compound containing solution is coated on the SiO.sub.2 film 4 by rotational coating. Then, solvent in the solution is removed by heating process. As a result, the surface of the SiO.sub.2 film 4 is covered by a SOG (Spin On Glass) film 5 which is made of silicon compound and has an even surface. Subsequently, the SOG film 5 is etched back so as to be left partially on concave portions of the SiO.sub.2 film 4. Both the SiO.sub.2 film 4 and the SOG film 5 serve as the interlayer insulating film so that an even surface of the interlayer insulating film can be derived. Then, the SiO.sub.2 film (not shown) is formed as a second layer by CVD process. A upper wiring (not shown) is formed on the second layer SiO.sub.2 film.\nIn the meanwhile, since the insulating film 4 formed between of the wiring 3, 3 serves as a dielectric film of a capacitor, wiring capacitance (parasitic capacitance) exists between two of the wiring 3.\nThe wiring capacitance C can be obtained according to following equation (1). The wiring capacitance C is proportional to a relative dielectric constant of the insulating film 4 located between two of the wiring 3, and is inversely proportional to an interval d between the wiring 3. In general, the relative dielectric constant of the SiO.sub.2 is about 4 to 5. In the equation (1), .epsilon..sub.0 denotes dielectric constant in vacuum, .epsilon..sub.r denotes relative dielectric constant, and S denotes an area of a side surface of the wiring 3. EQU C=.epsilon..sub.0 .epsilon.rS/d (1)\nThe interval d between the wiring 3, 3 becomes short and a multilayered structure of the wiring is advanced as high integration of the semiconductor integrated circuit device proceeds. As a result, the wiring capacitance is increased and thus an operation of the semiconductor device is further delayed."}
{"text": "Photovoltaic cells are widely used for electricity generation with one or more photovoltaic cells typically sealed within and interconnected in a module. Multiple modules may be arranged into photovoltaic arrays used to convert solar energy into electricity by the photovoltaic effect. Arrays can be installed on building rooftops and are used to provide electricity to the buildings and to the general grid."}
{"text": "The present invention relates to cosmetic and dermatological preparations, in particular those of the water-in-oil type, to processes for their preparation and to their use for cosmetic and medicinal purposes.\nThe human skin is man\"\"s largest organ and performs a number of vital functions. Having an average area of about 2 m2 in adults, it has a prominent role as a protective and sensory organ. The purpose of this organ is to transmit and avert mechanical, thermal, actinic, chemical and biological stimuli. In addition, it has an important role as a regulatory and target organ in human metabolism.\nThe main aim of skin care in the cosmetics sense is to strengthen or rebuild the skin\"\"s natural function as a barrier against environmental influences (e.g. dirt, chemicals, microorganisms) and against the loss of endogenous substances (e.g. water, natural fats, electrolytes), and also to assist its horny layer in its natural regeneration ability where damage has occurred.\nIf the barrier properties of the skin are impaired, increased resorption of toxic or allergenic substances or infection by microorganisms may result, leading to toxic or allergic skin reactions.\nAnother aim of skin care is to compensate for the loss by the skin of sebum and water caused by daily washing. This is particularly important if the natural regeneration ability is inadequate. Furthermore, skin care products should protect against environmental influences, in particular against sun and wind, and delay skin ageing.\nMedicinal topical compositions usually comprise one or more medicaments in an effective concentration. For the sake of simplicity, in order to clearly distinguish between cosmetic and medicinal use and corresponding products, reference is made to the legal provisions in the Federal Republic of Germany (e.g. Cosmetics Directive, Foods and Drugs Act).\nEmulsions are generally taken to mean heterogeneous systems which consist of two liquids which are immiscible or miscible with one another only to a limited extent, which are usually referred to as phases. In an emulsion, one of the two liquids is dispersed in the form of very fine droplets in the other liquid.\nIf the two liquids are water and oil and oil droplets are very finely dispersed in water, this is an oil-in-water emulsion (O/W emulsion, e.g. milk). The basic character of an O/W emulsion is determined by the water. In the case of a water-in-oil emulsion (W/O emulsion, e.g. butter), the principle is reversed, the basic structure being determined here by the oil.\nThe person skilled in the art is of course aware of a large number of ways to formulate stable W/O preparations for cosmetic or dermatological use, for example in the form of creams and ointments which can be spread in the range from room temperature to skin temperature, or as lotions and milks, which are more likely flowable in this temperature range. However, there are only a few formulations in the prior art which are of sufficiently low-viscosity that they would, for example, be sprayable.\nIn addition, low-viscosity preparations of the prior art frequently have the disadvantage that they are unstable, and are limited to a narrow field of application or a limited choice of feed materials. Low-viscosity products in which, for example, strongly polar oilsxe2x80x94such as the plant oils otherwise frequently used in commercially available productsxe2x80x94are sufficiently stabilized are therefore currently not on the market.\nThe term xe2x80x9cviscosityxe2x80x9d means the property of a liquid to resist the mutual laminar displacement of two neighbouring layers (internal friction). This so-called dynamic viscosity is nowadays defined according to xcex7=t/D as the ratio of shear stress to the velocity gradient perpendicular to the direction of flow. For Newtonian liquids, xcex7 is a material constant having the SI unit Pascal second (Paxc2x7s) at a given temperature.\nThe quotient xcexd=xcex7/xcfx81 from the dynamic viscosity xcex7 and the density xcfx81 of the liquid is referred to as the kinematic viscosity xcexd and is given in the SI unit m2/s.\nFluidity (xcfx86) is the inverse of viscosity (xcfx86=1/xcex7). In the case of ointments and the like, the use value is inter alia codetermined by the so-called tack. The tack of an ointment or ointment base or the like means its property to draw threads of varying lengths when a small sample is removed; accordingly, a distinction is made between short- and long-stretch substances.\nWhilst the graphical representation of the flow behaviour of Newtonian liquids at a given temperature produces a straight line, in the case of so-called non-Newtonian liquids considerable deviations often arise, depending on the velocity gradient D (shear rate xcex3) or the shear stress xcfx84. In these cases, the so-called apparent viscosity can be determined which, although not obeying the Newtonian equation, can be used to determine the true viscosity values by graphical methods.\nFalling-body viscometry is suitable only for investigating Newtonian liquids and gases. It is based on Stokes\"\"s law, according to which, for the falling of sphere through a liquid which flows around it, the dynamic viscosity xcex7 can be determined from   η  =            2      ⁢                                    r            2                    ⁢                      (                                          ρ                K                            -                              ρ                                  F                  ⁢                                      xe2x80x83                                    ⁢                  I                                                      )                          ·        g                    9      ·      v      \nwhere\nr=radius of the sphere, v=fall velocity, xcfx81K=density of the sphere, xcfx81FI=density of the liquid and g=acceleration of the fall.\nW/O emulsions with a high water content and any desired viscosity which moreover have storage stability, as is required for marketable products, in particular those with a relatively high water content of more than 75% by weight water content and nevertheless having very good sensory properties continue to be a weighty problem. As a result of the very high water content of the emulsions, they xe2x80x9cbreakxe2x80x9d on the skin particularly rapidlyxe2x80x94sensorily unpleasantxe2x80x94into their main constituents (hydrophilic and hydrophobic components). Solutions have hitherto only been attempted for W/O emulsions comprising predominantly nonpolar lipids. Accordingly, the supply of formulations of this type is extremely low.\nAn object of the present invention was to provide preparations which can be formulated in virtually any viscosities and which do not have the disadvantages of the prior art.\nAnother object of the present invention was to provide preparations which can be charged with a high content of water-soluble and/or water-miscible substances having cosmetic or dermatological activity, without impairing the galenical quality or other properties of the preparations.\nAccording to K. J. Lissant: The Geometry of High-Internal-Phase-Ratio Emulsions; Journal of Colloid and Interface Science 22, 462-468 (1966), emulsions with an internal phase of more than 70% are defined as so-called high internal phase emulsions. The preparation of stable, flowable water-in-oil emulsions with a water content of more than 80% is very difficult. In particular, xe2x80x9chigh internal phasexe2x80x9d W/O emulsions with a very high water content of more than 85% (xe2x80x9cvery high internal phasexe2x80x9d W/O emulsions) are not accessible.\nThe technique of varying the phase volume ratio (i.e. incorporating higher amounts of liquid lipids) which is usually used for water-in-oil emulsions can, because of the low lipid content, be used only to a limited extent in the case of high internal phase W/O emulsions, or not at all in the case of very high internal phase W/O emulsions.\nSurprisingly, it has been found that water-in-oil emulsions\n(a) with a content of water and optionally water-soluble substances totalling at least 75% by weight, and with a content of lipids, emulsifiers and lipophilic constituents totalling at most 25%, preferably at most 20%, based in each case on the total weight of the preparations,\n(b) whose oil phase is chosen from the group of lipids or lipid mixtures, where the total polarity of the lipid phase is between 20 and 30 mN/m,\n(c) comprising at least one interface-active substance, chosen from the group of alkylmethicone copolyols and/or alkyldimethicone copolyols,\n(d) optionally, but particularly when the content of water and water-soluble (hydrophilic) constituents is between 75 and 80% by weight, comprising one or more cationic polymers,\novercome the disadvantages of the prior art.\nIt is possible and advantageous to choose the total content of water and water-soluble substances of the W/O emulsions according to the invention to be greater than 80% by weight, in particular 85% by weight, in each case based on the total weight of the preparations.\nOne example of interface-active substances which are to be used particularly advantageously for the purposes of the present invention is cetyldimethicone copolyol, which is sold by Th. Goldschmidt AG under the trade name ABIL(copyright) EM 90.\nFurthermore, the emulsifier laurylmethicone copolyol has been found to be particularly advantageous, which is obtainable under the trade name Dow Corning(copyright) 5200 Formulation Aid from Dow Corning Ltd.\nThe total amount of the interface-active substances used according to the invention in the finished cosmetic or dermatological preparations is advantageously chosen from the range 0.1-30% by weight, preferably 0.25-5.0% by weight, in particular 0.75-3.5% by weight, based on the total weight of the preparations.\nAlthough it is known that W/O emulsions with a high water content can be produced using emulsifiers of the type described previously, the known prior art was nevertheless unable to indicate the route to the present invention.\nFor the purposes of the present disclosure, a general term for fats, oils, waxes and the like which is sometimes used is the term xe2x80x9clipidsxe2x80x9d, with which the person skilled in the art is entirely familiar. The terms xe2x80x9coil phasexe2x80x9d and xe2x80x9clipid phasexe2x80x9d are also used synonymously.\nOils and fats differ from one another in their polarity, which is difficult to define. It has already been proposed to adopt the interfacial tension towards water as a measure of the polarity index of an oil or of an oil phase. This means that the lower the interfacial tension between this oil phase and water, the greater the polarity of the oil phase in question. According to the invention, the interfacial tension is regarded as one possible measure of the polarity of a given oil component.\nThe interfacial tension is the force which acts on an imaginary line one meter in length in the interface between two phases. The physical unit for this interfacial tension is conventionally calculated from the force/length relationship and is usually expressed in mN/m (millinewtons divided by meters). It has a positive sign if it endeavours to reduce the interface. In the converse case, it has a negative sign.\nTable 1 below lists moderately polar lipids which are advantageous according to the invention as individual substances or else as mixtures with one another. The interfacial tensions towards water concerned are given in the last column. It is, however, also advantageous to use mixtures of higher and lower polarity and the like, provided it is ensured that the total polarity of the oil phase is in the claimed range.\nAccording to the teaching presented herewith, W/O emulsions are obtainable whose viscosity at 25xc2x0 C. is less than 5000 mPaxc2x7s (=millipascal seconds), in particular less than 4000 mPaxc2x7s, preferably less than 3500 mPaxc2x7s (HAAKE Viscotester VT-02).\nFor the purposes of the present invention, the oil phase can additionallyxe2x80x94provided the features listed in the patent claims are consideredxe2x80x94advantageously comprise substances chosen from the group of esters of saturated and/or unsaturated, branched and/or unbranched alkanecarboxylic acids having a chain length of from 3 to 30 carbon atoms and saturated and/or unsaturated, branched and/or unbranched alcohols having a chain length of from 3 to 30 carbon atoms and from the group of esters of aromatic carboxylic acids and saturated and/or unsaturated, branched and/or unbranched alcohols having a chain length of from 3 to 30 carbon atoms. Such ester oils can then advantageously be chosen from the group consisting of isopropyl myristate, isopropyl palmitate, isopropyl stearate, isopropyl oleate, n-butyl stearate, n-hexyl laurate, n-decyl oleate, isooctyl stearate, isononyl stearate, isononyl isononanoate, 2-ethylhexyl palmitate, 2-ethylhexyl laurate, 2-hexyldecyl stearate, 2-octyldodecyl palmitate, oleyl oleate, oleyl erucate, erucyl oleate, erucyl erucate and also synthetic, semisynthetic and natural mixtures of such esters, such as, for example, jojoba oil.\nThe oil phase can also be chosen advantageously from the group of branched and unbranched hydrocarbons and hydrocarbon waxes, silicone oils, dialkyl ethers, the group of saturated or unsaturated, branched or unbranched alcohols, and fatty acid triglycerides, namely the triglycerol esters of saturated and/or unsaturated, branched and/or unbranched alkanecarboxylic acids having a chain length of from 8 to 24 carbon atoms, in particular 12-18 carbon atoms. The fatty acid triglycerides can, for example, be advantageously chosen from the group of synthetic, semisynthetic and natural oils, e.g. olive oil, sunflower oil, soya oil, peanut oil, rapeseed oil, almond oil, palm oil, coconut oil, palm kernel oil and the like.\nIf desired fatty and/or wax components which are to be used in the oil phasexe2x80x94as secondary constituents in a minor amountxe2x80x94can be chosen from the group of vegetable waxes, animal waxes, mineral waxes and petrochemical waxes. Examples which are favourable according to the invention are candelilla wax, carnauba wax, japan wax, esparto grass wax, cork wax, guaruma wax, rice germ oil wax, sugarcane wax, berry wax, ouricury wax, montan wax, jojoba wax, shea butter, beeswax, shellac wax, spermaceti, lanolin (wool wax), uropygial grease, ceresin, ozocerite (earth wax), paraffin waxes and microcrystalline waxes.\nOther advantageous fatty and/or wax components are chemically modified waxes and synthetic waxes such as, for example, those obtainable under the trade names Syncrowax HRC (glyceryl tribehenate), Syncrowax HGLC (C16-36-fatty acid triglyceride) and Syncrowax AW 1C (C18-36 fatty acid) from CRODA GmbH, and also montan ester waxes, Sasol waxes, hydrogenated jojoba waxes, synthetic or modified beeswaxes (e.g. dimethicone copolyol beeswax and/or C30-50-alkyl beeswax), polyalkylene waxes, polyethylene glycol waxes, but also chemically modified fats, such as, for example, hydrogenated vegetable oils (for example hydrogenated castor oil and/or hydrogenated coconut fatty glycerides), triglycerides, such as, for example, trihydroxystearin, fatty acids, fatty acid esters and glycol esters, such as, for example, C20-40-alkyl stearate, C20-40-alkylhydroxystearoyl stearate and/or glycol montanate. Also advantageous are certain organosilicon compounds, which have similar physical properties to the specified fatty and/or wax components, such as, for example, stearoxytrimethylsilane.\nIf desired, the fatty and/or wax components can be present either individually or as a mixture.\nAny desired mixtures of such oil and wax components can also be used advantageously for the purposes of the present invention. In some instances, it can also be advantageous to use waxes, for example cetyl palmitate, as the lipid component of the oil phase.\nOf the hydrocarbons, paraffin oil, hydrogenated polyolefins (e.g. hydrogenated polyisobutene), squalane and squalene can be used advantageously for the purposes of the present invention.\nAccording to the invention, emulsions which are particularly advantageous are those which are characterized in that the oil phase consists of at least 50% by weight, preferably of more than 75% by weight, of at least one substance chosen from the group consisting of vaseline (petrolatum), paraffin oil and polyolefins, and of the latter, preference is given to polydecenes.\nThe oil phase can advantageously additionally have a content of cyclic or linear silicone oils or consist entirely of such oils, although it is preferable to use an additional content of other oil phase components apart from the silicone oil or the silicone oils.\nCyclomethicone (octamethylcyclotetrasiloxane) can be used advantageously. However, other silicone oils can also be used advantageously for the purposes of the present invention, for example hexamethylcyclotrisiloxane, polydimethylsiloxane and poly(methylphenylsiloxane).\nSurprisingly, it has been found, in particular, that by adding from 0.01 to 10%, preferably 0.25-1.25%, of suitable cationic polymers, it is possible to prepare stable, flowable xe2x80x9cvery high internal phase emulsionsxe2x80x9d which have excellent sensory properties.\nSuitable cationic polymers are, for example, cationic cellulose derivatives (e.g. Polymer JR 400(copyright) from Amerchol), cationic starch, copolymers of diallylammonium salts and acrylamides, quaternized vinylpyrrolidone/vinylimidazole polymers (e.g. Luviquat(copyright) from BASF), condensation products of polyglycols and amines, quaternized collagen polypeptides (e.g. Lamequat(copyright) L from Grxc3xcnau-Henkel), quaternized wheat polypeptides, polyethyleneimine, cationic silicone polymers, copolymers of adipic acid with dimethylaminohydroxypropyldiethylenetriamine, copolymers of acrylic acid with dimethyldiallylammonium chloride (e.g. Merquat(copyright) 550 from Chemviron), polyaminopolyamides, cationic chitin derivatives, cationic guar gum (e.g. Jaguar(copyright) CBS from Hoechst Celanese), quaternized ammonium salt polymers (e.g. Mirapol(copyright) AD-1 from Miranol), and cationic biopolymers, such as, for example, chitosan (average molecular weight from 50,000 to 2,000,000 g/mol [determined by means of gel permeation chromatography] and a degree of deacylation of from 10 to 99% [determined by means of 1H-NMR]).\nThe aqueous phase of the preparations according to the invention in some instances advantageously comprises alcohols, diols or polyols of low carbon number, and ethers thereof, preferably ethanol, isopropanol, propylene glycol, glycerol, ethylene glycol, ethylene glycol monoethyl ether or monobutyl ether, propylene glycol monomethyl, monoethyl or monobutyl ethers, diethylene glycol monomethyl or monoethyl ethers and analogous products, and also alcohols of low carbon number, e.g. ethanol, isopropanol, 1,2-propanediol and glycerol.\nA particular advantage of the present invention is that it permits high concentrations of polyols, in particular glycerol, to be used.\nParticularly advantageous preparations are also obtained when antioxidants are used as additives or active ingredients. According to the invention, the preparations advantageously comprise one or more antioxidants. Antioxidants which are favourable but which are nevertheless optional may be all antioxidants which are customary or suitable for cosmetic and/or dermatological application.\nThe antioxidants are advantageously selected from the group consisting of amino acids (e.g. glycine, histidine, tyrosine, tryptophan) and their derivatives, imidazoles, (e.g. urocanic acid) and their derivatives, peptides, such as D,L-carnosine, D-carnosine, L-carnosine and their derivatives (e.g. anserine), carotenoids, carotenes (e.g. xcex1-carotene, xcex2-carotene, \"psgr\"-lycopene) and their derivatives, lipoic acid and its derivatives (e.g. dihydrolipoic acid), aurothio-glucose, propylthiouracil and other thiols (e.g. thioredoxin, glutathione, cysteine, cystine, cystamine and their glycosyl, N-acetyl, methyl, ethyl, propyl, amyl, butyl and lauryl, palmitoyl, oleyl, xcex3-linoleyl, cholesteryl and glyceryl esters) and their salts, dilauryl thiodipropionate, distearyl thiodipropionate, thiodipropionic acid and its derivatives (esters, ethers, peptides, lipids, nucleotides, nucleosides and salts) and sulphoximine compounds (e.g. buthionine sulphoximines, homocysteine sulphoximine, buthionine sulphones, penta-, hexa-, heptathionine sulphoximines) in very low tolerated doses (e.g. pmol to xcexcmol/kg), and also (metal) chelating agents (e.g. xcex1-hydroxy fatty acids, palmitic acid, phytic acid, lactoferrin), xcex1-hydroxy acids (e.g. citric acid, lactic acid, malic acid), humic acid, bile acid, bile extracts, bilirubin, biliverdin, EDTA, EGTA and their derivatives, unsaturated fatty acids and their derivatives (e.g. xcex3-linolenic acid, linoleic acid, oleic acid), folic acid and its derivatives, ubiquinone and ubiquinol and their derivatives, vitamin C and derivatives (e.g. ascorbyl palmitate, Mg ascorbyl phosphate, ascorbyl acetate), tocopherols and derivatives (e.g. vitamin E acetate), vitamin A and derivatives (vitamin A palmitate) and coniferyl benzoate of benzoin resin, rutinic acid and its derivatives, ferulic acid and its derivatives, butylated hydroxytoluene, butylated hydroxyanisole, nordihydroguaiacic acid, nordihydroguaiareticic acid, trihydroxybutyrophenone, uric acid and its derivatives, mannose and its derivatives, zinc and its derivatives (e.g. ZnO, ZnSO4), selenium and its derivatives (e.g. selenomethionine), stilbenes and their derivatives (e.g. stilbene oxide, trans-stilbene oxide), and the derivatives (salts, esters, ethers, sugars, nucleotides, nucleosides, peptides and lipids) of said active substances which are suitable according to the invention.\nFor the purposes of the present invention, oil-soluble antioxidants can be used particularly advantageously.\nA surprising property of the present invention is that preparations according to the invention are very good vehicles for cosmetic or dermatological active ingredients into the skin, preferred active ingredients being antioxidants which are able to protect the skin against oxidative stress. Preferred antioxidants are vitamin E and its derivatives and vitamin A and its derivatives.\nThe amount of antioxidants (one or more compounds) in the preparations is preferably from 0.001 to 30% by weight, particularly preferably 0.05-20% by weight, in particular 0.1-10% by weight, based on the total weight of the preparation.\nIf vitamin E and/or its derivatives are used as the antioxidant or antioxidants, their respective concentrations are advantageously chosen from the range of 0.001-10% by weight, based on the total weight of the formulation.\nIf vitamin A or vitamin A derivatives or carotenes or their derivatives are used as the antioxidant or antioxidants, their respective concentrations are advantageously chosen from the range of 0.001-10% by weight, based on the total weight of the formulation.\nThe person skilled in the art is of course aware that cosmetic preparations are in most cases inconceivable without the customary auxiliaries and additives. The cosmetic and dermatological preparations according to the invention can, accordingly, also comprise cosmetic auxiliaries, as are customarily used in such preparations, for example bodying agents, stabilizers, fillers, preservatives, perfumes, antifoams, dyes, pigments which have a colouring action, thickeners, surface-active substances, emulsifiers, emollients, moisturizers and/or humectants, anti-inflammatory substances, additional active ingredients such as vitamins or proteins, sunscreens, insect repellants, bactericides, virucides, water, salts, antimicrobial, proteolytic or keratolytic substances, medicaments or other customary constituents of a cosmetic or dermatological formulation such as alcohols, polyols, polymers, foam stabilizers, organic solvents or also electrolytes.\nThe latter can be chosen, for example, from the group of salts containing the following anions: chlorides, also inorganic oxo element anions, of these, in particular sulphates, carbonates, phosphates, borates and aluminates. Electrolytes based on organic anions are also advantageous, e.g. lactates, acetates, benzoates, propionates, tartrates, citrates, amino acids, ethylenediamine-tetraacetic acid and salts thereof and others. Preferred cations of the salts are ammonium, alkylammonium, alkali metal, alkaline earth metal, magnesium, iron or zinc ions. It does not need to be mentioned that only physiologically acceptable electrolytes should be used in cosmetics. Particular preference is given to potassium chloride, sodium chloride, magnesium sulphate, zinc sulphate and mixtures thereof.\nCorresponding requirements apply mutatis mutandis to the formulation of medicinal preparations.\nThe W/O emulsions according to the invention can be used as a basis for cosmetic or dermatological formulations. The latter can have the customary composition and be used, for example, for the treatment and care of the skin and/or the hair, as lip care product, as deodorant product and as make-up or make-up remover product in decorative cosmetics or as a sunscreen preparation. For use, the cosmetic and dermatological preparations according to the invention are applied to the skin and/or the hair in a sufficient amount in a manner customary for cosmetics or dermatological compositions.\nFor the purposes of the present invention, cosmetic or topical dermatological compositions can accordingly, depending on their composition, be used, for example, as a skin protection cream, cleansing milk, sunscreen lotion, nourishing cream, day or night cream, etc. In some circumstances it is possible and advantageous to use the compositions according to the invention as a base for pharmaceutical formulations.\nThe low-viscosity cosmetic or dermatological compositions according to the invention can, for example, be in the form of preparations which can be sprayed from aerosol containers, squeezable bottles or by means of a pump device, or in the form of a liquid composition which can be applied by means of roll-on devices, but also in the form of an emulsion which can be applied from normal bottles and containers.\nSuitable propellants for cosmetic or dermatological preparations which can be sprayed from aerosol containers for the purposes of the present invention are the customary known readily volatile, liquefied propellants, for example hydrocarbons (propane, butane, isobutane), which can be used alone or in a mixture with one another. Compressed air is also used advantageously.\nThe person skilled in the art is of course aware that there are propellants which are non-toxic per se which would be suitable in principle for realizing the present invention in the form of aerosol preparations, but which must nevertheless be avoided because of their unacceptable impact on the environment or other accompanying circumstances, in particular fluorinated hydrocarbons and chlorofluorocarbons (CFCs).\nCosmetic and dermatological preparations which are in the form of a sunscreen are also favourable. As well as the active ingredient combinations according to the invention, these preferably additionally comprise at least one UV-A filter substance and/or at least one UV-B filter substance and/or at least one inorganic pigment.\nFor the purposes of the present invention, however, it is also advantageous to provide cosmetic and dermatological preparations whose main purpose is not protection against sunlight, but which nevertheless have a content of UV protectants. Thus, for example, UV-A or UV-B filter substances are usually incorporated into day creams.\nUV protectants, like antioxidants and, if desired, preservatives, also effectively protect the preparations themselves against decay.\nPreparations according to the invention can advantageously comprise further substances which absorb UV radiation in the UV-B range, the total amount of filter substances being, for example, from 0.1% by weight to 30% by weight, preferably from 0.5 to 10% by weight, in particular from 1.0 to 6.0% by weight, based on the total weight of the preparations, in order to provide cosmetic preparations which protect the hair and/or the skin from the whole region of ultraviolet radiation. They can also be used as sunscreens for the hair or the skin.\nIf the emulsions according to the invention contain UV-B filter substances, the latter may be oil-soluble or water-soluble. Examples of oil-soluble UV-B filters which are advantageous according to the invention are:\n3-benzylidenecamphor derivatives, preferably 3-(4-methylbenzylidene)camphor, 3-benzylidenecamphor;\n4-aminobenzoic acid derivatives, preferably 2-ethylhexyl 4-(dimethylamino)benzoate, amyl 4-(dimethylamino)benzoate;\nesters of cinnamic acid, preferably 2-ethylhexyl 4-methoxycinnamate, isopentyl 4-methoxycinnamate;\nesters of salicylic acid, preferably 2-ethylhexyl salicylate, 4-isopropylbenzyl salicylate, homomenthyl salicylate;\nderivatives of benzophenone, preferably 2-hydroxy-4-methoxybenzophenone, 2-hydroxy-4-methoxy-4xe2x80x2-methylbenzophenone, 2,2xe2x80x2-dihydroxy-4-methoxybenzophenone;\nesters of benzalmalonic acid, preferably di(2-ethylhexyl) 4-methoxybenzalmalonate;\nderivatives of 1,3,5-triazine, preferably 2,4,6-trianilino(p-carbo-2xe2x80x2-ethyl-1xe2x80x2-hexyloxy)-1,3,5-triazine.\nThe list of said UV-B filters, which may be used in combination with the active ingredient combinations according to the invention is of course not intended to be limiting.\nIt can also be advantageous to formulate lipodispersions according to the invention with UV-A filters which have hitherto been customarily present in cosmetic preparations. These substances are preferably derivatives of dibenzoylmethane, in particular 1-(4xe2x80x2-tert-butylphenyl3-(4xe2x80x2-methoxyphenyl)propane-1,3-dione and 1-phenyl-3-(4xe2x80x2-isopropylphenyl)propane-1,3-dione.\nCosmetic and dermatological preparations according to the invention can also comprise inorganic pigments which are customarily used in cosmetics for protecting the skin against UV rays. These are oxides of titanium, zinc, iron, zirconium, silicon, manganese, aluminium, cerium and mixtures thereof, and modifications in which the oxides are the active agents. Particular preference is given to pigments based on titanium dioxide.\nFurther constituents which can be used are:\nfats, waxes and other natural and synthetic fatty substances, preferably esters of fatty acids with alcohols of low carbon number, e.g. with isopropanol, propylene glycol or glycerol, or esters of fatty alcohols with alkanoic acids of low carbon number or with fatty acids;\nalcohols, diols or polyols of low carbon number, and their ethers, preferably ethanol, isopropanol, propylene glycol, glycerol, ethylene glycol, ethylene glycol monoethyl or monobutyl ethers, propylene glycol monomethyl, monoethyl or monobutyl ethers, diethylene glycol monomethyl or monoethyl ethers and analogous products.\nPreparations according to the invention can also comprise active ingredients (one or more compounds) which are chosen from the group: acetylsalicylic acid, atropine, azulene, hydrocortisone and derivatives thereof, e.g. hydrocortisone-17 valerate, vitamins, e.g. ascorbic acid and derivatives thereof, vitamins of the B and D series, very favourably vitamin B1, vitamin B12 and vitamin D1, but also bisabolol, unsaturated fatty acids, namely the essential fatty acids (often also called vitamin F), in particular xcex3-linolenic acid, oleic acid, eicosapentanoic acid, docosahexanoic acid and derivatives thereof, chloramphenicol, caffeine, prostaglandins, thymol, camphor, extracts or other products of a vegetable and animal origin, e.g. evening primrose oil, star flower oil or currant seed oil, fish oils, cod-liver oil or also ceramides or ceramide-like compounds etc. It is also advantageous to choose the active ingredients from the group of refatting substances, for example purcellin oil, Eucerit(copyright) and Neocerit(copyright).\nThe amount of such active ingredients (one or more compounds) in the preparations according to the invention is preferably from 0.001 to 30% by weight, particularly preferably 0.05-20% by weight, in particular 1-10% by weight, based on the total weight of the preparation.\nThe examples below serve to illustrate the present invention without limiting it. The numerical values in the examples refer to percentages by weight, based on the total weight of the respective preparations."}
{"text": "An electronically controlled brake system for a motor vehicle is known, which electrically controls the braking forces of braking devices for the respective wheels, namely, the hydraulic pressures supplied to wheels cylinders for driving the braking devices, in accordance with the amount of a braking operation, e.g., the amount of depression of a brake pedal. An example of a control apparatus for this type of brake system is disclosed in Japanese Laid-open Patent Publication No. 2004-243983.\nIn the vehicular brake control system as disclosed in the above-identified publication, when the vehicle operator or driver operates the brake pedal, a master cylinder generates a hydraulic pressure in accordance with the amount of operation of the brake pedal, and a part of the hydraulic oil flows into a stroke simulator so that the amount of operation of the brake pedal (represented by, e.g., pedal stroke) is controlled in accordance with the pedal pressure applied to the brake pedal (i.e., the force with which the brake pedal is depressed by the driver). In the meantime, the brake control system sets a target deceleration of the vehicle in accordance with the pedal stroke detected by brake ECU, and determines the distribution of the braking forces applied to the respective wheels, based on which certain hydraulic pressures are applied to the respective wheel cylinders.\nIn the known vehicular brake control system as described above, the master cylinder for generating a hydraulic pressure commensurate with the amount of operation of the brake pedal is equipped with the stroke simulator into which a part of the hydraulic oil is fed so as to control the amount of operation of the brake pedal. The master cylinder is also provided with pressurizing mechanisms for pressurizing the hydraulic oil to be supplied to the wheel cylinders of the four wheels via master cut valves. Since the pressurizing mechanisms are provided for the wheel cylinders of the respective wheels, the overall hydraulic system is complicated in construction, resulting in an increase in the manufacturing cost."}
{"text": "Technical Field\nThe present disclosure relates to metal organic frameworks comprising zinc (II) ions and second metal ions, such as iron (II) ions, cobalt (II) ions, and copper (II) ions as nodes or clusters and coordinated 1,3,5-benzenetricarboxylic acid struts or linkers between them. Additionally, the present disclosure relates to processes for producing the metal organic frameworks and their application as catalysts in methods for the oxidation of cyclic hydrocarbons, such as toluene, cyclohexane, and methylcyclohexane.\nDescription of the Related Art\nThe “background” description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description which may not otherwise qualify as prior art at the time of filing, are neither expressly or impliedly admitted as prior art against the present invention.\nMetal organic frameworks (MOFs) are crystalline, highly porous coordination polymers which are comprised of inorganic units (i.e. metals) coordinated to rigid organic fragments. They were built by a node linker approach that was first described by Robson [Hoskins, B. F.; Robson, R. J. Am. Chem. Soc. 1990, 112, 1546.—incorporated herein by reference in its entirety]. This method uses metal ions as nodes and organic ligands as linkers. In this case, a metal ion with a preferred coordination number and geometry combines with divergent ligand molecules to form an extended 1D, 2D or 3D network or networks. The metal linker interactions vary widely and have included ion-ion interactions, ion-dipole, dipole-dipole, hydrogen bonding, anion-π interactions, π-π interactions, as well as van-der Waals interactions [Burrows, A. D.; Chan, C. W.; Chowdhry, M. M.; McGrady, J. E.; Mingos, D. M. P. Chem. Soc. Rev. 1995, 329; and Classens, C. G.; Stoddart, J. F. J. Phys. Org. Chem. 1997, 10, 254.—each incorporated herein by reference in its entirety]. The strength of these interactions has been proven to directly influence the overall stability of the resulting framework [Braga, D. Chem. Comm. 2003, 2751-2754.—incorporated herein by reference in its entirety]. Extended networks can be developed from the basic principles guiding the formation of coordination complexes. FIG. 1 shows a representation of the formation of these coordination networks. Hence, the choice of metal center affects the resulting framework structure because a given metal has preference for a specific geometry and coordination environment.\nA second approach to the synthesis of MOFs was described by Yaghi [Eddaoudi, M.; Li, H.; Reineke, T.; Fehr, M.; Kelley, D.; Groy, T. L.; Yaghi, O. M. Topics in Catalysis 1999, 9, 105.; and Eddaoudi, M.; Moler, D. B.; Li, H.; Chen, B.; Reineke, T. M.; O'Keeffe, M.; Yaghi, O. M. Acc. Chem. Res. 2001, 34, 319.—each incorporated herein by reference in its entirety]. This involves the use of multiple organic ligands and linkers and metal ions or clusters (secondary building units, SBUs) as nodes. FIG. 2 shows a schematic illustration of this MOF synthesis and structure. Most often, MOFs are easily synthesized by means of hydrothermal or solvothermal synthesis which involve high temperature self-assembly in a high boiling organic solvent or water in closed vessels [Yaghi, O. M.; O'Keeffe, M.; Ockwig, N. W.; Chae, H. K.; Eddaoudi, M.; Kim, J. M Nature 2003, 423, 705-714.; and Stock, N.; Biswas, S. Chem. Rev., 2012, 112, 933-969.; and Yaghi, O. M.; Li, H. L. J. Am. Chem. Soc., 1995, 117, 10401-10402.; and Lin, W. B.; Wang, Z. Y.; Ma, L. J. Am. Chem. Soc., 1999, 121, 11249-11250.; and Chui, S. S. Y.; Lo, S. M. F.; Charmant, J. P. H.; Orpen, A. G.; Williams, I. D. Science, 1999, 283, 1148-1150.—each incorporated herein by reference in its entirety]. However, these methods typically require long reaction times, from several hours up to several days, depending upon the nature of the ligand, the reaction solvent, reagent concentrations and reaction temperature. They can also be produced using the microwave assisted process which allows the large scale synthesis of MOFs in a few minutes [Feldblyum, J.; Liu, M.; Gidley, D.; Matzger, A. J. Am. Chem. Soc. 2011, 133, 18257-18263.; and Dybtsev, D. N.; Nuzhdin, A. L.; Chun, H.; Bryliakov, K. P.; Talsi, E. P.; Fedin, V. P.; Kim, K. Angew. Chem., Int. Ed., 2006, 45, 916-920.; and Ni, Z.; Masel, R. I.; J. Am. Chem. Soc., 2006, 128, 12394-12395.—each incorporated herein by reference in its entirety].\nMOF-5 is one of the first series of MOFs to be reported and fully characterized. It consists of an octahedral secondary building unit (SBU) which is made from Zn4O (CO2)6 as an inorganic unit which is comprised of four ZnO4 tetrahedra with a common vertex and six carboxylate groups. These octahedral SBUs are joined together by benzene linkers [Chalati, T.; Horcajada, P.; Gref, R.; Couvreur, P.; Serre, C. J. Mater. Chem., 2011, 21, 2220-2227.—incorporated herein by reference in its entirety]. These unique units lead to a perfect cubic network whose vertices comprise the SBUs and the edges of which are made up of the benzene linkers. FIG. 3A shows the structure of MOF-5 in an extended 3D cubic framework. This compound was synthesized from Zn(II) and 1,4-benzenedicarboxylic acid (BDC) under organic conditions predetermined to form the SBU in situ. FIG. 3B shows the topology of the MOF structure as a ball-and-stick model. FIG. 3C shows the structure represented by (OZn4)O12 clusters (tetrahedrons) joined by BDC ions (connectors). Since the benzene links and the SBUs appear to be relatively rigid and large entities, the resulting structure has exceptionally high porosity (as indicated by its sorption) and stability [Yaghi, O. M., Sun, Z., Richardson, D. A. & Groy, T. L J. Am. Chem. Soc. 1994, 116, 802-808.—incorporated herein by reference in its entirety].\nMOFs have been found to possess unique properties like high surface areas (up to 10400 m2/g) and tunable pores that can be used in various potential applications such as gas storage, catalysis, separation, and drug delivery [Li, H., Eddaoudi, M., O'Keeffe, M. & Yaghi, O. M Nature 1999, 402, 276-279.; and Silva, P.; Valente, A. A.; Rocha, J.; Paz, F. A. A. Cryst. Growth Des., 2010, 10, 2025-2028.; and Wang, Z.; Chen, G.; Ding, K. Chem. Rev., 2009, 109, 322.; and Corma, A.; Garci'a, H.; Xamena, F. X. Chem. Rev., 2010, 23, 1126.; and Li, J. R.; Kuppler, R. J.; Zhou, H. C. Chem. Soc. Rev., 2009, 38, 1477.—each incorporated herein by reference in its entirety]. They have also been found to be candidates for other applications like microelectronics, sensing, optics, micromotors, molecular rotors, and bioreactors [Horcajada, P.; Chalati, T.; Serre, C.; Gillet, B.; Sebrie, C.; Baati, T.; Eubank, J. F.; Heurtaux, D.; Clayette, P.; Kreuz, C.; Chang, J. S.; Hwang, Y. K.; Marsaud, V.; Bories, P. N.; Cynober, L.; Gil, S.; Fe'rey, G.; Couvreur, P.; Gref, R. Nat. Mater., 2010, 9, 172-178.; and Talin, A. A.; Centrone, A.; Ford, A. C.; Foster, M. E.; Stavila, V.; Haney, P.; Kinney, R. A.; Szalai, V.; Gabaly, F. E.; Yoon, H. P.; Le'onard, F.; Allendorf, M. D. Science, 2014, 343, 66.; and Kreno, L. E.; Leong, K.; Farha, O. K.; Allendorf, M.; Van Duyne, R. P.; Hupp, T. Chem. Rev., 2012, 112, 1105.; and Cui, Y.; Yue, Y.; Qian, G.; Chen, B. Chem. Rev., 2012,112, 1126.; and Ikezoe, Y.; Washino, G.; Uemura, T.; Kitagawa, S.; Matsui, H. Nat. Mater., 2012, 11, 1081.; and Comotti, A.; Bracco, S.; Ben, T.; Qiu S.; Sozzani, P. Angew. Chem., Int. Ed., 2014, 53, 6655.—each incorporated herein by reference in its entirety]. Presently, the most highly recognized applications of MOFs, however, have been in the areas of gas storage and separation. The strong interest in this research area stems from the urgent need to develop viable technologies for hydrogen fuel storage for commercial use as well as to control the concentration of CO2 in the atmosphere.\nHigh surface area and the possibility of varied structural modification amongst other desirable physical and chemical properties make it possible for MOFs to efficiently catalyze a broad range of reactions. The MOFs are usually modified by a method that is broadly known as post-synthetic-modification (PSM). PSM makes it possible to incorporate a highly diverse range of different functional groups making it largely free of the restrictions resulting from the synthetic conditions of the MOFs. PSM also allows the introduction of multiple metal ions into a single framework in a combinatorial manner, enabling an effective way to systematically fine tune and optimize MOF properties [Doherty, C. M.; Grenci, G.; Riccoo', R.; Mardel, J. I.; Reboul, J.; Furukawa, S.; Kitagawa, S.; Hill, A. J.; Falcaro, P. Adv. Mater., 2013, 25, 4701.—incorporated herein by reference in its entirety]. The process in which new metal sites are incorporated into a MOF framework is known as transmetallation or post-synthetic metal exchange and the MOFs produced by this method can be described as isostructural MOFs with similar structural frameworks but different metal ions. FIG. 4 illustrates a general scheme for the post-synthetic modification of MOFs. This synthetic method can be used to obtain certain MOFs that cannot be obtained via conventional synthetic methods. Cation exchange also helps to enhance the properties of some MOFs by making it possible to incorporate a more useful metal site thereby improving some of their physical and chemical properties hence giving them more interesting applications. For example, HKUST-1 has a surface area of about 1500 m2/g and it contains Cu2+ but the isostructure can be made by substituting the Cu2+ with other metals. A unique property of transmetallation lies in the fact that new MOFs can be obtained by complete or partial substitution of metal ions within the framework without altering the morphology of the MOFs. FIG. 5 shows a schematic representation for post synthetic metal exchange or transmetallation. This process serves as an alternative, typically milder route for accessing new MOFs when conventional synthesis at high temperature fails [Wang, Z.; Cohen, S Chem. Soc. Rev., 2009, 38, 1315-1329.—incorporated herein by reference in its entirety]. This substitution occurs at the metal nodes, often called the inorganic clusters or secondary building units (SBUs). Although the metal ions are integral parts of the MOFs' structures, they can be replaced either completely or partially within hours or days without necessarily affecting the MOFs' structures [Dinca, M.; Long, J. R. J. Am. Chem. Soc., 2007, 129, 11172-11176.—incorporated herein by reference in its entirety]. Transmetallation changes the properties of the MOFs and also makes them useful for other important applications especially in catalysis. The transmetallated MOFs have multiple properties having separate metal sites that can be utilized for specific catalytic conversions of organic molecules.\nThe development of efficient new catalysts is still a serious challenge in chemical research. Hence the increasing demand for safer and energy saving reaction routes promotes the need to develop new materials towards the global aim of combating serious environmental challenges that stem from several industrial processes. Hence, catalyst development is an ever growing area of research. Recently, chemists have endorsed MOFs as viable heterogeneous catalysts to channel the course of new and existing chemical reactions to reduce industrial wastes and enable greener chemical processes. This work facilitates better understanding of physical and chemical processes such as surface interactions and facilitates novel concepts and ideas for the next generation of catalysts. The role of heterogeneous catalysts either in chemical or petrochemical industries cannot be overemphasized. They reduce the enormous wastes that are associated with homogeneous catalysts and also reduce cost due to their reusability. These heterogeneous catalysts occur in a different phase from the substrates and predominantly work base on an adsorption mechanism. The heterogeneous catalysts are mostly solids on which liquid or gaseous reaction mixtures are adsorbed. The active site may be either a planar exposed metal surface, a crystal edge with imperfect metal valence or a complicated combination of the two. Thus, not only most of the volume, but also most of the surface of a heterogeneous catalyst may be catalytically inactive. The dependence of catalytic activity on surface area and pore volume makes MOF catalysts a viable area of chemical research. Investigating the nature of the active sites requires technically challenging research. Thus, studies relating to new metal and ligand combinations for catalysis continue.\nThese highly porous, crystalline MOFs have some of the catalytically important properties of zeolites like uniform cavity and pore sizes as well as medium to large internal surface areas [Wang, L, J.; Deng, H.; Furukawa, H.; Gandara, F.; Cordova, K. E.; Peri, D.; Yaghi, O. M. Inorg. Chem. 2014, 53, 5881-5883.—incorporated herein by reference in its entirety]. Unlike zeolites, vast chemical varieties of MOFs can be synthesized due to the presence of infinite organic linkers. This suggests that the catalytic niche of MOFs is likely to be high value added reactions such as production of specific enantiomers, sensitive molecules, as well as production of fine chemicals, which require specific and tunable catalytic sites [Furukawa, H.; Ko, N.; Go, Y. B.; Aratani, N.; Choi, S. B.; Choi,E.; Yazaydin, A.; Snurr, R. Q.; O'Keeffe, M.; Kim, J.; Yaghi, O. M. Science, 2010, 329, 424.; and Pan, L.; Adams, K. M.; Hernandez, H. E.; Wang, X.; Zheng, C.;Hattori, Y.; Kaneko, K. J. Am. Chem. Soc. 2003,125, 3062.—each incorporated herein by reference in its entirety]. Despite the various interesting and compelling recent developments in MOF catalysis, the area of MOF catalysis is still in an immature phase. Many researchers have likened MOF catalysis to enzyme catalysis, aiming towards the development of catalytic chemistry in the direction of an “artificial enzyme”. Overall, the uniqueness of MOFs over other materials is yet to be fully illustrated since they have been reported to be of use as catalysts in the chemical or petrochemical industries.\nOxidation reactions are among the most important chemical conversions in industries and laboratories. Conversion of abundant and cheap hydrocarbons like toluene, cycloalkanes and methylcyclohexane into more valuable chemicals like aldehydes, ketones and acids stands as a significant process for consideration [Hayashi, H.; Cote, A. P.; Furukawa, H.; O'Keeffe, M.; Yaghi, O. M. Nat. Mater. 2007, 6,501.—incorporated herein by reference in its entirety]. Among these useful transformations, the direct oxidation of toluene to produce benzaldehyde is an attractive process. Toluene oxidation gives a mixture of oxygenated products like benzoic acid, benzyl alcohol, benzaldehyde and cresols. Commercially, benzaldehyde is mainly produced by the chlorination of toluene followed by the hydrolysis process, which generates large amounts of toxic acidic/basic discard solutions, leading to equipment corrosion and environmental pollution. Furthermore, the benzaldehyde produced by this route is not qualified to synthesize some high quality compounds such as perfumes or pharmaceuticals because the product contains chlorine [Friedrich Brihne and Elaine Wright “Benzaldehyde” in Ullmann's Encyclopedia of Industrial Chemistry, 2002, Wiley-VCH, Weinheim.—incorporated herein by reference in its entirety]. Therefore, there is a clear need to develop new materials to improve the selective oxidation of toluene as an alternative route to produce benzaldehyde and consequently benzoic acids.\nAlkanes are naturally abundant and cheap carbon containing raw material which serve as attractive substrates for the production of value added organic chemicals (alcohols, ketones, aldehydes and carboxylic acids) [B. Retcher, J. S. Costa, J. Tang, R. Hage, P. Gamez, J. Reedijk, J. Mol. Catal. A. 2008, 286, 1-5.—incorporated herein by reference in its entirety]. Unfortunately, the chemical inertness of these compounds is a considerable limitation towards their vast application for direct syntheses of oxygenated products under relatively mild conditions. However, a proper metal catalyst and an appropriate oxidizing agent, as well as properly controlled reaction conditions, can lead to the development of a cleaner and more efficient chemical industry. Today, over a billion tons of cyclohexanone and cyclohexanol are produced every year and they are mostly used for the synthesis of Nylon-66 and Nylon-6 [M. Musser “Cyclohexanol and Cyclohexanone” in Ullmann's Encyclopedia of Industrial Chemistry, Wiley-VCH, Weinheim, 2005.—incorporated herein by reference in its entirety].\nThe structure of HKUST-1 and Zn-HKUST-1 metal organic frameworks has been investigated. Chui and coworkers were the first to synthesize HKUST-1 at the Hong Kong University of Science and Technology. HKUST-1 was discovered to be a highly porous metal coordination polymer [Cu3(BTC)2(H2O)3]n (where BTC is benzene-1,3,5-tricarboxylate) which has interconnected [Cu2(O2CR)4] units (where R is an aromatic ring), C18H12O15Cu3. It creates a 3-dimensional system of channels with a pore size of approximately 1 nm and an accessible porosity of about 40% in the solid. The single crystal structural analysis of HKUST-1 revealed that the polymer framework is composed of dimeric cupric tetracarboxylate building units, with a Cu—Cu inter-nuclear separation of 2.628(2) Å . FIG. 6 shows the crystal structure of the dicopper (II) tetracarboxylate building block containing two axial aqua ligands. The framework was found to be neutral because the twelve carboxylate oxygens from the two BTC ligands bind to four coordination sites for each of the three Cu2+ ions of the formula unit. Hence each Cu atom completes is pseudo-octahedral coordination sphere with the presence of axial aqua ligands opposite to the Cu—Cu dimer. The tetracarboxylate unit provides a structural motif with potential four-fold symmetry, and the trimesic acid provides a three-fold symmetry element. FIG. 7 shows a HKUST-1 secondary building unit (SBU) demonstrating the tbo net topology and the paddlewheel structure. The origin of the nanochannels can be clearly considered to arise from the formation of larger octahedral secondary building units (SBUs). The main SBU in HKUST-1 is the octahedral unit with Cu2 at its 6 vertices and 4 trimesate ions tetrahedrally disposed as “panels” for four of the eight triangular faces of the octahedron. FIG. 8 shows the [Cu3(TMA)2(H2O)3] unit viewed along the cell body diagonal, demonstrating a hexagonal-n shaped window at the intersection of the nanopores. FIG. 9 shows the polymer framework and nanochannels with four-fold symmetry.\nThe Zn-HKUST-1 that contains Zn2+ ions is analogous to HKUST-1. Analysis of Zn-HKUST-1 by powder X-ray diffraction and gas sorption shows the retention of crystalline structure but negative nitrogen uptake at 77 K due to a dense surface layer that prevents the passage of small molecular species into the crystal framework [Bhunia, M. K, Hughes, J. T. Fettinger, J. C. and Navrotsky, A. Langmuir 2013, 29, 8140-8145.—incorporated herein by reference in its entirety]. The prevalence of zinc paddlewheels in a variety of MOFs, such as the previously discussed HKUST-1 suggests that Zn I a promising metal with which an isostructural analog to Cu-HKUST-1 can be constructed.\nThe application of metal organic frameworks has been investigated. Over the past decades, transition metal complexes comprising mainly phosphine ligands, salen or salophen ligands, pincer ligands and N-heterocyclic carbenes (NHCs) have had a remarkable impact on catalysis. Examples include the Nobel prize winning Noyori asymmetric hydrogenation, Sharpless oxidations (Sharpless epoxidation, Sharpless asymmetric dihydroxylation, and Sharpless oxyamination) as well as Jacobsen epoxidation [Noyori, R Adv. Synth. Catul 2003, 345, 12, 15-41.; and Katsuki, T.; Sharpless, K. B J Am. Chem. Soc., 1980, 102, 5974-5976.; and Jacobsen, E. N.; Marko, T.; Mungall, W. S.; Schroeder, G.; Sharpless, K. B. J. Am. Chem. Soc., 1988, 110 1968-1970.; and Sharpless, K. B; Patrick, D. W.; Truesdale, L. K.; Biller, S. A. J. Am. Chem. Soc., 1975, 97, 2305 2307.; and Zhang, W.; Loebach, J. L.; Wilson, S. R.; Jacobsen, E. N. J. Am. Chem. Soc., 1990, 112, 7, 2801-2803.—each incorporated herein by reference in its entirety]. Most of these metal complexes have only been successful as homogeneous catalysts which have significant disadvantages like the difficulty of product separation, poor reusability and toxicity. Relatively insoluble and stable materials like zeolites, metals and metal oxides are widely used as heterogeneous catalyst on the industrial scale. Notable conversions like the Harber-Bosch, Contact process, Ostwald process, steam reforming, petrochemical reactions Ziegler Natta polymerization make use of these insoluble materials. In terms of easy post reaction separation, these materials have largely been successful. However, there is still an urgent need for the development of new materials that will be more energy efficient, more tunable, and more environmentally friendly.\nMOFs have shown great catalytic prospects for a wide range of reactions due to the diversity in their structures, low toxicity, reusability and cost effectiveness. In fact MOFs are among the best candidates in bridging the gap between homogeneous and heterogeneous catalysis. The metals in the MOF structure often act as Lewis acids especially when the frameworks are activated by removing the coordinated labile solvent molecules or counter ions [Han, J. W.; Hill, C. L.; J. Am. Chem. Soc., 2007, 129, 15094.—incorporated herein by reference in its entirety]. Fujita, et al. first reported the catalytic activity of a 2D Cd(II) based MOF for the cyanosilylation of aldehydes. Thy obtained the unsaturated metal cluster by removing two water molecules from the octahedral structure of [Cd(4,4′-bpy)2(H2O)2].(NO3)2.4H2O [Fujita, M.; Kwon, Y. J.; Washizu, S.; Ogura, K. J. Am. Chem. Soc. 1994, 116, 1151-1152.—incorporated herein by reference in its entirety]. Fe(BTC) has also been used as a heterogeneous catalyst for the selective methylation of primary aromatic amines using dimethyl carbonate, efficient oxidation of benzylic compounds using t-butyl-hydroperoxide as oxidizing agent, and oxidation of thiols to disulfides [Dhakshinamoorthy, A.; Alvaro, M.; Garcia, H. Appl. Catal. A: General 2010, 378, 19-25.; and Dhakshinamoorthy, A.; Alvaro, M.; Garcia, H. J. Catal. 2009, 267, 1-4.; and Dhakshinamoorthy, A.; Alvaro, M.; Garcia, H. Chem. Commun., 2010, 46, 6476-6478.—each incorporated herein by reference in its entirety]. Seo, et al. first reported asymmetric catalysis using a homochiral MOF, [Zn(μ3-O)(1-H)6.2H3O.12H2O] for transesterification reactions. It was also the first MOF demonstrating that the organic linker embedded into a pore can catalyze an asymmetric reaction [Seo, J. S.; Whang, D.; Lee, H.; Jun, S. I.; Oh, J.; Jeon, Y. J.; Kim, K. Nature, 2000, 404, 982-986.—incorporated herein by reference in its entirety]. Lin, et al. also reported the activity of a homochiral non-interpenetrating MOF which was constructed in finite 1-dimensional [Cd(μ-Cl)2]n zigzag chains with axial bipyridine bridging ligands containing orthogonal secondary functional groups [Wu, C.; Hu, A.; Zhang, L.; Lin, W. J. Am. Chem. Soc., 2005, 127, 8940-8941.—incorporated herein by reference in its entirety]. The chiral secondary functional groups were used to generate a heterogeneous asymmetric catalyst for the addition of diethyl zinc to aromatic aldehydes to afford chiral secondary alcohols at up to 93% enantiomeric excess (ee).\nMore specifically, the catalytic activity of HKUST-1 has been investigated. HKUST-1 has been particularly well recognized for its high catalytic activity especially when the axial aqua ligands are removed via activation. FIG. 10 illustrates schematically the activation of HKUST-1. Activation gives unsaturated metal sites without affecting the rigid framework of the MOF [Schlichte, K.; Kratzke, T.; Kaskel, S.; Microporous and Mesoporous Materials. 2004, 73: 81-85.; and Lien T. L. Nguyen, Tung T. Nguyen, Khoa D. Nguyen, Nam T. S. Phan Applied Catalysis A: General, 2012, 425, 44-52.—each incorporated herein by reference in its entirety]. Schlichte, et al. first reported the catalytic activity of HKUST-1 when they used the HKUST-1 MOF for the trimethylcyanosilation of benzaldehyde. The open framework of this MOF was activated by removing the two water molecules from axial positions in the octahedral framework. The activated MOF afforded up to 57% conversion of benzaldehyde reaching a selectivity of 89% at 313 K. Nguyen, et al. studied the activity of HKUST-1 for the aza-michael reaction in which amines were reacted with α,β-unsaturated carbonyl groups to prepare β-amino carbonyl compounds and their derivatives. They achieved excellent conversions up to 100% under relatively mild conditions in the presence of 5 mol % activated catalyst. Fourier transform infrared spectroscopy (FT-IR) and powder X-ray diffraction (PXRD) analysis revealed that the catalyst could be reused several times without a significant reduction in its catalytic potency. Atomic absorption spectroscopy showed that the reaction was not influenced by homogeneous catalysis resulting from leached active species.\nUsing the same material, Phan, et al. was able to react phenols and aryl iodides to form diaryl ethers in an Ullman-type coupling reaction. The heterogeneous reaction leads to high conversion using 5 mol % catalyst in the presence of MeONa as a base. The used catalyst was facilely recovered from the reaction mixture using simple filtration and could be reused without significant degradation [Nam T. S. Phan, Tung T. Nguyen, Chi V. Nguyen, Thao T. Nguyen Applied Catalysis A: General., 2013, 457, 69-77.—incorporated herein by reference in its entirety]. Phan, et al. also reported the highly efficient activity of HKUST-1 for the C-arylation of acetylacetone in the presence of aryl iodides to obtain aryl ketones as major products. HKUST-1 was confirmed to be a true heterogeneous catalyst as there was no effect of homogeneous catalysis of active species leaching into the reaction mixture [Nam T. S. Phan, Tung T. Nguyen, Phuong Ho, and Khoa D. Nguyen ChemCatChem. 2013, 5, 1822-1831.—incorporated herein by reference in its entirety]. Dang, et al. studied the catalytic activity of HKUST-1 for the synthesis of propargylamine via direct oxidative C—C coupling reaction using C—H functionalization between phenylacetylene and N,N-dimethylaniline to give N-methyl-N-(3-phenylprop-2-ynyl)benzenamine as the principal product [Giao H. Dang, Duy T. Nguyen, Dung T. Le, Thanh Truong, Nam T. S. Phan Journal of Molecular Catalysis A, 2014, 300, 306.—incorporated herein by reference in its entirety]. The copper catalyzed reaction afforded 96% conversion after 180 minutes at 120° C. in the presence of 5 mol % copper-based catalyst. The used catalyst was recovered from the reaction medium by filtration and reused for the coupling reaction. Similarly, HKUST-1 has been used as a catalyst for the direct oxidative amination of sp2 C—H bonds. The reaction involves the use of N-methylmorpholine oxide (NMO) as oxidizing agent in the presence of primary or secondary amine as coupling pairs with DMF as solvent at 90-100° C. [Nga T. T. Tran, Quan H. Tran, Thanh T. Journal of Catalysis, 2014, 320, 9-15.—incorporated herein by reference in its entirety].\nIn view of the forgoing, one object of the present disclosure is to provide relatively cheap and environmentally friendly metal organic framework catalysts designed towards laboratory and industrial scale catalytic applications as opposed to uses in gas storage and carbon capture. Metal organic framework catalysts provide great potential for new and existing chemical reactions for shortest route organic conversions that reduce industrial wastes and enable greener chemical processes through their high tunability, high surface area, stability, and reusability. Specifically, this disclosure is focused on isostructural HKUST-1 metal organic frameworks comprising zinc (II) metal ions and second metal ions, such as iron (II) ions, cobalt (II) ions, and copper (II) ions linked by 1,3,5-benzentricarboxylic acid to form a porous coordination network as polyhedral crystals. This disclosure provides the transmetallation preparation of these metal organic frameworks by post-synthetic metallic exchange of the solvothermally synthesized Zn-HKUST-1 metal organic framework. An additional aspect of the present disclosure is application of these metal organic frameworks as catalysts in methods for the liquid phase oxidation of cyclic hydrocarbons, such as, toluene, cyclohexane, and methylcyclohexane. It is envisioned that the metal organic frameworks of the present disclosure will exhibit strong activity in terms of conversion of the cyclic hydrocarbon, selectivity for a desired oxidized cyclic hydrocarbon product, and reusability. Overall, the metal organic frameworks are envisaged to exhibit strong potential utility as catalysts that aid the increasing demand for safer and energy saving reaction routes that additionally reduce and minimize the adverse environmental impact of industrial wastes."}
{"text": "This invention relates to interactive television program guide systems, and more particularly, to interactive television program guide systems that allow users to record programs and program guide data on a media server.\nCable, satellite, and broadcast television systems provide viewers with a large number of television channels. Viewers have traditionally consulted printed television program schedules to determine the programs being broadcast at a particular time. More recently, interactive television program guides have been developed that allow television program information to be displayed on a user's television. Interactive television program guides allow the user to navigate through television program listings using a remote control. In a typical program guide, various groups of television program listings are displayed in predefined or user-defined categories. Listings are typically displayed in a list, grid, or table.\nProgram listings and other program guide data are typically provided by a satellite uplink facility to a number of cable system headends. Each headend distributes the program guide data to a number of users. Interactive television program guides are typically implemented on the users' set-top boxes. A typical set-top box is connected to a user's television and videocassette recorder. Program guide systems that allow users to record programs are described, for example, in. Ellis et al. U.S. patent application Ser. No. 08/924,239, filed Sep. 5, 1997, which is hereby incorporated by reference herein in its entirety.\nSuch systems are deficient in a number of respects. The processing and storage capabilities of the set-top box and videocassette recorder are generally limited. This, in turn, may limit the functionality of the guide. Videocassette recorders may add significantly to the cost of in-home television equipment. They are mechanical systems, prone to failure, and their proper operation relies in part on users setting them properly (i.e., remembering to put a tape in). Setting up a VCR to work cooperatively with a program guide may be a complex process in some systems and may frustrate users. In addition, users cannot record multiple programs simultaneously without having multiple VCRs, and recording one program while watching another typically requires additional hardware or an enhanced set-top box. Systems in which program guides allow users to record one program while watching another are described, for example, in Lemmons et al. U.S. patent application Ser. No. 60/089,487, filed Jun. 16, 1998 which is hereby incorporated by reference herein in its entirety.\nSystems that use hard disk technology to store programs have also been developed. Interactive television program guide systems that have digital storage devices are described, for example, in Hassell et al. U.S. patent application Ser. No. 09/157,256, filed Sep. 17, 1998. Hard-disk based products have also been developed by TiVo, Inc. of Sunnyvale, Calif., and Replay Networks, Inc. of Mountainview, Calif. These systems are deficient in that they require additional hardware in the user's home that may significantly increase the cost of the user's home television equipment. Such systems also do not allow users to record multiple programs simultaneously without having multiple devices in the home.\nSome current television platforms support a return path between the set-top boxes and the headends. Client-server based program guides have been developed in which set-top boxes act as clients that communicate with servers located at headends via return paths. The servers typically provide program listings information (e.g., program titles and broadcast times) to the set-top boxes in response to requests that are generated by the set-top boxes.\nVideo-on-demand (VOD) systems have also been developed. Such systems typically record all programs that are distributed by a headend, or only a chosen subset of programs. In the former approach, vast amounts of storage is required at the server to ensure that all possible videos desired by users will be available. In the latter approach, users are limited to viewing only those programs that the headend operator decided to record.\nIn one cable television system known as “The Box,” viewers can call a cable operator to request the airing of a particular music video from an archive of music videos. The requested videos are broadcast on a dedicated channel for all viewers to see. This approach suffers from a number of deficiencies. First, users do not have the convenience of ordering programs using their television equipment. In addition, viewers must wait for the airing of other viewers' selections before theirs are aired. Also, viewers are limited to selecting videos from only those archived. Some viewers may desire music videos that are not shared in the archive. These viewers have no way of selecting a video for archiving or viewing other videos.\nIn a cable television system known as “Your Choice TV,” television programs that have previously aired are repeated on a dedicated set of television channels. Viewers may order a repeat of a program for a small fee. However, the repeats are aired at scheduled times, which may not be convenient for the viewer. Also, desired programs may not be available to viewers because the cable system operator decides which programs to record.\nIt is therefore an object of the present invention to provide a program guide system that allows users to direct a server to record certain programs that later may be played back to the user on demand."}
{"text": "This invention relates to a system for determining the operating characteristics of a rod-pumped oil well and making automatic control decisions based on those determinations. Previous methods for detection and control of a condition known as \"pump off\" have evaluated data on a dynagraph, which displays measured polished rod load and measured or calculated polished rod position. Some prior systems, such as U.S. Pat. No. 4,286,925 to Standish, test for pump-off by determining whether the load at a particular point in the downstroke exceeds a preset or user-adjustable limit. Other systems have measured the area within the dynagraph for one full stroke (called a card), which represents work done by the pump, and compared that area against a limit or a \"test card.\" U.S. Pat. No. 3,951,209 to Gibbs, for example, discloses a method of integrating the entire area within a dynagraph. U.S. Pat. No. 4,015,469 to Womack discloses a method of integrating equal portions of the upstroke and downstroke. End Devices, Inc., in its device known as the Model 107DC, disclosed a method of integrating the lower half of the dynagraph, i.e., the downstroke. Our U.S. Pat. No. 4,583,915 discloses a method of integrating a portion of the area below the dynagraph, i.e., the space between measured load and minimum load during a selected time period.\nEach of those methods share the shortcomings that they are difficult to adjust and sometimes falsely detect pump-off when the well is in fact full. For example, when a well is shut down for a long period of time, such as for service work, the fluid level may rise in the annuls. That rise in fluid level reduces the hydrostatic head required to lift the fluid to the surface, just as if the well were shallower. When the pump is restarted, therefore, it needs to do less work, and the area inside the dynagraph may be reduced to the point that pump-off will be detected, even though the pump is full. U.S. Pat. No. 4,302,157, issued to Welton, et al., teaches a method for preventing false pump-off detection in which pump-off is recognized only after the pump first operates normally with a full pump. However, the method in that patent is complex and uses a simplistic \"rule of thumb,\" rather than a highly discriminating test. All of the patents cited above are incorporated herein by reference.\nThe present invention provides an improved system for controlling a well and detecting pump-off."}
{"text": "Terminals come in a wide variety of configurations for different environments. Telecommunications terminals were optimized for use with copper drop wires. Through an evolutionary process, a plurality of different wire sizes and kinds of wire are employed to connect the subscriber to the phone company. The gauges can vary from 12-30 AWG gauge copper based wires some of which may have steel cores. Most common wire sizes are between 16-26 gauge. Standard conventional terminals have a threaded binding post embedded in a suitable dielectric base material. The drop wire is stripped of its insulation, formed in a \"C\" shape, and connected to the terminal post by nuts and washers. This procedure is craft sensitive and time consuming. These terminals, while extremely low cost, suffer high maintenance and repair costs due to corrosion, e.g., oxidation on the exposed wires and binding posts. In addition, wet/humid weather, or periodic water submersion, or salt/fog corrosion, or dew on insect nests can cause cross talk or signal loss. A serviceman is dispatched in response to a customer complaint only to find that the problem has disappeared due to the evaporation of the moisture.\nInsulation displacement terminals, which cut through the wire insulation without requiring wire stripping, address the time consuming installation problems. However, many of these terminals also suffer from corrosion, often more severe than standard binding post and washer/nut terminals. In addition, overnotching of the conductor may sever or severely weaken the drop wire rendering it subject to premature failure. The premature failure of the drop wires significantly increases the telephone company's repair costs. In addition, initial use of the terminal on a large drop wire may damage the terminal and preclude its subsequent use on smaller drop wires.\nU.S. application Ser. Nos. 07/462,173 filed Jan. 8, 1990, now U.S. Pat. No. 5,069,636 07/231,755 filed Aug. 13, 1988, now abandoned, and U.S. Pat. No. 4,846,721 solved many of these problems with a uniquely formed electrical connection and sealing system utilizing a reenterable gel material. The gel sealing system utilized in the terminals are, inter alia, disclosed in U.S. Pat. Nos. 4,865,905; 4,864,725; 4,600,261; 4,634,207; 4,643,924; and 4,690,831. However, these blocks were only readily adaptable to wires, e.g., drop wires. These preceding specifications and patents are completely incorporated herein by reference for all purposes.\nAlthough these terminals can form electrical contacts to a plurality of different wire sizes, an even easier, more versatile system would be desirable. In addition, with an increasing use of computers, and the overcrowding of conventional telephone lines limiting the capacity of the copper plant system architecture, a terminal which can be subsequently adapted after installation, to include circuit protection, filters, cable TV, optical fibers, and/or digitizing of the telephone signal, e.g., digital added main lines, (DAML), to increase the number of different phone numbers available through each twisted pair of wires, would be highly desirable.\nIt would also be desirable to have a modular drop wire unit which is simple for the craftsperson to install and keyed to the block to avoid misinsertions. It would additionally be desirable to have a unit which can adequately seal in a reversible fashion between the drop wire module and the block itself for subsequent additions to the block such as circuit protection, and the like.\nThe inability to provide a secure weatherproof system, less subject to outside plant failures, has forced the use of many add-on features, such as fused circuit protection and DAML, at the customer's protected base location. This renders them subject to tampering. The weatherproof system could be centrally located underground, on a telephone pole, or in a centrally located outside enclosure, e.g., a pedestal enclosure. As the system is upgraded, the new features can be plugged into the terminal through the common plug connection.\nBy a weatherproof system, we mean a modular system that meets the requirements of a present day standard outside nut and washer or insulation displacement terminal that can operate for prolonged periods in the outside plant environment. The block of the invention alone without connections or a block in combination with modules is capable of exhibiting resistance to the elements even after repeated connections/disconnections during its lifetime. At least five connections/disconnections and most likely ten or more can be expected during a terminal's lifetime. Suitable tests are ISO 846 (ASTM G-21) (1985) for fungi and bacteria. The dielectric strength of the electrical insulation in the block must be capable of passing IEC 243 (ASTM D-149) (1981). The block must also must be able to withstand corrosive effects of materials and the effects of liquid chemicals meeting the requirements of ASTM D-2671 (1985) and ISO 175 (ASTM D-543) (1987). While in use, the contact resistance cannot vary greater than the parameters permitted by ASTM B 539-80 (1985 revised) as well as withstanding salt fog testing according to ASTM B-117 (1985). Finally, to be effective the block must be able to operate while exposed to sunlight and water while remaining able to meet the standards of ASTM G-53-84 (1988) and ASTM D-257 (1983).\nThus, it would be highly desirable to have a weatherproof modular system capable of providing repeated reentry of the plug for the drop wire, fuse, digitally added mainline (DAML) apparatus, and the like, without subjecting the terminal to either corrosion or electrical failure.\nThe process of making the terminals is labor intensive and time consuming. The process requires the individual termination of the copper wires within the cable to the back of the terminal and thereafter potting. An embodiment of the present invention provides for the cable wires to be attached to the terminal posts without precutting. In the prior art, if the wires are cut to an improper length either a splice or a total rewiring of the terminal block was usually required. To provide a lower cost terminal and a more reliable and less labor intensive manufacturing process, it would be highly desirable to have a block which can be tested during the wiring process and prior to wire cutting and assembly."}
{"text": "Cheesecake is one type of food which has maintained its popularity amongst the baker and consumer alike. Cheesecake is typically made by the home baker by mixing a variety of ingredients including a major portion of milk or milk products to form a batter, placing the batter into a baking container, placing the baking container into an oven, sometimes within a water bath, at a preset temperature for a certain period of time sufficient to bake the cheesecake, removing the cheesecake from the oven, allowing the cheesecake to cool to room temperature in an ambient environment and placing the cheesecake in a cold environment for at least twenty-four hours before serving.\nHome-baked cheesecake is usually consumed within 24 to 72 hours of the time it is removed from the oven as it cannot be stored for a very long time, even in a cold environment such as the one provided by a refrigerator since many changes occur rather quickly to the palatability and texture of the cheesecake. Such changes may even have serious consequences on the health of those who consume the cheesecake if the time for which the cheesecake has been stored in the refrigerator or the temperature of the refrigerator have been such that they have allowed for bacteria and mold to develop and grow either on or within the body of the cheesecake.\nAs the demand for cheesecake has grown steadily with the growing population and as more and more home bakers' time is taken up with working outside the home, commercial bakers have sought to mass-produce cheesecake. Such mass production and attempts at marketing this delectable dessert has led to the expansion of equipment and adoption of techniques which include: (a) automated on-line mixing of ingredients; (b) specialized baking forms and molds; (c) computer-controlled industrial size ovens; (d) automated transfer of ingredients throughout the baking method through the use of conveyors and other movers; and (e) automated packaging and labeling of baked goods. Despite all of these and other features, which have increased the speed by which commercial cheesecake can be made, such features often have a negative impact on the taste and quality of the cheesecake produced. More importantly, however, such techniques and equipment have failed to change the process of mixing, baking, cooling, serving the cheesecake and/or prolonging the storage time and/or shelf life of the cheesecake prior to consumption.\nCommercial cheesecake, i.e. factory baked cheesecake generally sold in retail bakeries and supermarkets, has a very short shelf-life due to its high susceptibility to spoilage. Such spoilage is probably due to its milk and other highly perishable ingredients. Thus, commercially baked cheesecake must either be refrigerated, frozen, or baked with preservatives. Freezing most definitely spoils the taste. Refrigeration must be carefully monitored and is effective for only very short periods of time. As for using chemical preservatives as ingredients for the cheesecake to prolong its shelf life, it is not a desirable alternative because (a) such preservatives affect the flavor and other physical characteristics of the cheesecake; and (b) in view of the general health conscious attitudes adopted by the consuming public at large, use of such preservatives will most likely lead to a drop in sales and consumability of the commercially baked cheesecake. More importantly, however, is that no matter which one of these three methods is used, the shelf or storage life of commercial cheesecake is relatively short and, when exceeded, results in a great deal of spoilage and return of cheesecakes to the baker for credit or replacement. Such returns translate to economic loss for the bakers and, by extension, to an increase of cheesecake prices to the consumer.\nIt is, therefore, an object of the present invention to provide a baking process and method for commercially available cheesecake which results in a product of such purity and lasting characteristics that it may be stored, transported and displayed without refrigeration.\nIt is a further object of the present invention to provide a commercially baked cheesecake having a long shelf life under normal atmospheric conditions, without any freezing, refrigeration or preservatives whatsoever.\nStill another object of the invention is to provide a baking process for commnercially available cheesecake, the cheesecake having a flavor, aroma and taste similar to that of freshly made home-baked cheesecake, without needing preservatives, freezing or refrigeration to maintain its relatively long shelf life.\nThe objects and advantages together with others will be apparent from the following disclosure of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a sheet-bundle conveying method, a sheet-bundle conveying apparatus used in such a sheet-bundle conveying method, and a bookbinding apparatus having such a sheet-bundle conveying apparatus. More particularly, it relates to a sheet-bundle conveying method, a sheet-bundle conveying apparatus used in such a sheet-bundle conveying method, and a bookbinding apparatus having such a sheet-bundle conveying apparatus, in which, after output sheets are book-bound, the book-bound sheet-bundle is conveyed.\n2. Related Background Art\nIn conventional bookbinding apparatuses for an image forming apparatus such as a printer, a sheet-bundle conveying means served to convey a sheet bundle comprised of imaged sheets from an alignment portion to a bookbinding portion and then from the bookbinding portion to a containing portion. In particular, a shiftable tray for conveying and containing the sheet-bundle to the containing portion was bulky for maintaining a large shifting amount of the tray and for ensuring the rigidity of the tray."}
{"text": "In order to understand the overall behavior of data traffic in a network domain, it is important to understand the characteristics of the individual data traffic flows in the network. The characteristics of the individual data traffic flows can include information about the individual data packets or metrics related to the set of packets in each data flow. These metrics can include measurements such as average data packet size, types of encapsulation of data packets, delay between sending and receiving data packets across the network domain, variations in throughput for data flows across the network (i.e., jitter), bandwidth utilization for a data flow, identification of a path for each data flow and similar metrics and data related to the each data flow. Once this sort of network analytic data has been collected for a data flow it can be analyzed in conjunction with similar data collected for the other data flows traversing the network to build an image or model of the behavior of data traffic across the network as a whole.\nCollecting network analytic data is performed at a subset of nodes within the network domain. These nodes can be selected or placed within the network to provide a coherent picture of the data traffic in the network domain. For example, the nodes that collect the network analytic data can be ingress or egress nodes in the network domain at the edge of network domain enabling end to analysis of data flows. Nodes in the interior can also be utilized to identify internal bottlenecks and similar issues within the network domain. However, collecting and reporting network analytic data can consume significant resources at the node that is performing the collection.\nThe addition of this task to a node can have a negative effect on its performance. Thus, it is important that the collection of network analytic data at each participating node be as efficient as possible. Many solutions for collecting network analytic data are inefficient and unworkable. A naïve implementation would be to examine each data packet that traverses a node in the network. However, this would be to computationally intensive and would severely impact the performance of the node. In many cases, sampling is performed at a node to select a representative subset of the data traffic of the data flow to utilize for network analytics. However, the use of sampling prevents the tracking of the same data packet across the network domain. This limits the types of analysis that can be performed, because the information about data packets will vary from node to node and analysis based on end to end tracking of individual data packets cannot be performed."}
{"text": "1. Field of the Invention\nA method for controlling an automatic stall-prevention of regulated-rotation-speed low-pressure centrifugal fan and a control automatics therefor.\n2. Description of Background Art\nNumerous compressors and centrifugal fans of different types, depending on their various applications, are known in the art. Particularly in industry, centrifugal fans, compressors and radial fans are widely used to achieve a pressure difference in piping. A problem for all centrifugal fan is generally known to be stalling. In other words, stalling is a characteristic state for all centrifugal fans, which occurs when the volume flow rate is too small in relation to the speed of rotation of the impeller. In this case the angle of incidence between the flow and the blade changes to be so disadvantageous that the flow disengages from the surface of the blade. Backflow is then able to occur in the blade passage and the impeller loses its pressure-increasing ability.\nIn this way, cyclical pressure fluctuations are generated that excite the natural frequencies of the structure surrounding, among other things, the piping. Pressure fluctuations create fatigue loading in the piping and in the structures. In addition, the temperature of the flow can significantly rise when the impeller, via the losses, continuously gives thermal energy to the gas, but the effective flow can be very small.\nProblems occur particularly in processes in which resistance to the flow greatly changes. When using a centrifugal fan for producing low pressure in these types of processes, a stalling state must be prevented by giving leakage air to the centrifugal fan in a controlled manner.\nConventionally, centrifugal fans that have a constant speed of rotation have been used. In this case stalling is prevented with an automatic leakage air valve, which receives control from the current of the drive motor of the centrifugal fan. In the stalling state the current of the drive motor is smaller than in the normal operating range. The electric current also fluctuates strongly. The control logic of the centrifugal fan can easily be programmed to detect a stalling state and to eliminate it by means of leakage air.\nAnother problem is that in a stalling state the low pressure also fluctuates with a rapid cycle, and this situation is detrimental from the viewpoint of the process. A device that has a regulated speed of rotation has been launched in the low-pressure centrifugal fan market as a new technology, the stall control of which cannot be implemented with conventional technology."}
{"text": "Polymer particles are used in a polymer matrix to manage light diffusion. For example, U.S. Pat. No. 7,547,736 describes the use of particles having an average particle size of 15 to 70 microns to provide a frosted look and textured surface, and U.S. Pat. No. 8,163,827 describes a high light transmission diffusion screen having inorganic pigments and refractive index-matched particles.\nPoint light sources create a visible shape of the light source, and there is often a desire to hide the light source shape, creating a more diffuse lighting. As used herein, “point light source” means any shaped source of electromagnetic radiation in the 3,600-7,700 Angstrom range. This includes, but is not limited to, incandescent, fluorescent, neon, argon and LED light sources.\nLight emitting diodes (LEDs) are being increasingly favored as a light source, since they use far less electricity and produce higher luminous efficacy than standard incandescent or fluorescent light bulbs. LEDs provide a very bright, point light source, yet the output (when correlated color temperature CCT is close to 7000° K) often appears harsh and causes an uncomfortable glare. This is a challenge for lighting designers, as many applications demand an illumination and low glare. Light devices that integrate one or more LEDs include, for example, luminaires for commercial or residential lighting, motor vehicle illumination devices (at the front or rear), indicating panels, luminous displays, spotlights, street lighting, box-letters, etc.\nThese luminous devices typically consist of a light source and a cover (also called a lens or a diffuser) made of a plastic whose function is to mask and protect the light source, while still ensuring good transmission of the light emitted by the light source. The plastic may be colored or may have decorative elements or patterns. The cover also has the function of scattering the emitted light so that the illumination is softened and not dazzling. The scattering of the light emitted by the light source can be achieved by dispersing scattering particles of organic or mineral nature in the plastic.\nReplacing a conventional light source with an LED results in a modification of the illumination. This is because an LED, especially an LED having a high luminous flux, exhibits directional illumination, whereas, for example the illumination of a neon tube is from 0 to 360°. In addition, the emission spectrum of an LED is completely different from that of a conventional light source.\nWO 2006/100126 describes a thermoplastic cover with dispersed beads for use with LEDs to form luminous devices. 3-30% of scattering particles are dispersed in a transparent plastic. The particles can be inorganic or organic and have mean diameters of from 0.5 to 100 microns. There is no description of combinations of particle size and loading, and no teaching of hiding power.\nThe addition of scattering particles helps to soften the effect of the LED light source, but the scattering also reduces light transmission. Some LED lighting cover manufacturers add pigments, such as BaSO4, SiO2, CaCO3, AL2O3, TiO2 and ZnO (U.S. Pat. No. 4,418,986), to the cover to increase the hiding power, though this can dramatically decrease the light transmission. BaSO4 and refractive-index-matched beads are used in a light diffusion (TV) screens in U.S. Pat. No. 8,163,827. Ground-up cell-cast sheet forms irregular particles with a wide particle size distribution, and unsatisfactory light transmission properties for LED diffusion sheets.\nWO 14/055330 describes an attempt to balance the light transmission and hiding power for an LED diffuser, using a combination of different sized plastic beads.\nThere is a need to provide excellent transmission, haze and diffusion properties with a lower particle loading.\nIt has now surprisingly been found that high efficiency diffusion coverings (film, sheet and profiles) for use with LED point light sources can be achieved by careful selection of the organic particle size, thickness of the covering, and difference in refractive index between the organic particles and the transparent polymer matrix. By matching the particle characteristics to achieve a simulated (calculated) transmittance, the particles can be used in an efficient manner, with lower loadings for a given haze, transmission and light diffusion.\nSelected properties of the high efficiency diffusion coverings can be further improved by the use of white pigment particles of a selected particle size and composition, and also by the use of optical brighteners."}
{"text": "Pressure transducers have been employed in a myriad of applications. One such transducer is the capacitance manometer which provides very precise and accurate measurements of pressure of a gas, vapor or other fluid. Applications include precise pressure measurement and high-precision gas and vapor delivery systems, which have become very important in many industrial applications, for example in the semiconductor industry for wafer and chip fabrication, although other applications are known. Such fluid delivery systems typically also include, but are not limited to, devices such as mass flow controllers (MFCs) and mass flow verifiers (MFVs) to regulate and/or monitor the flow of gases and vapors.\nCapacitance manometers typically use a flexible diaphragm forming or including an electrode structure and a fixed electrode structure spaced from the diaphragm so as to establish a capacitance there between. Variation in pressure on one side of the diaphragm relative to the pressure on the opposite side of the diaphragm causes the diaphragm to flex so that the capacitance between the electrode structure of the diaphragm and the fixed electrode structure varies as a function of this differential pressure. Usually, the gas or vapor on one side of the diaphragm is at the pressure being measured, while the gas or vapor on the opposite side of the diaphragm is at a known reference value, whether at atmosphere, or some fixed high or low (vacuum) pressure, so that the pressure on the measuring side of the diaphragm can be determined as a function of the capacitance measurement.\nWhen a capacitance manometer is used to measure pressure and the diaphragm flexes to provide a capacitance change, it is expected that when the pressure on the measuring side of the diaphragm returns to the same pressure as the reference side (the “zero state”), the instrument will indicate a “zero” reading. However, over time, for various reasons, the reading will drift from zero when the manometer returns to a zero state. Accordingly, the manometer reading must be zeroed and calibrated from time to time to return the reading to zero for the zero state. For prior art capacitance manometers, determination of whether the manometer needs to be zeroed (e.g., have its zero-reading calibrated) has typically been decided on a time related basis, e.g., every six months, or on a routine basis, e.g., every preventative maintenance cycle (PM), or never. The choice of how often to zero the manometer is typically left to the discretion of an operator or user (e.g., a semiconductor process engineer). Such a decision has typically been based on a human inspection of previous drifts.\nAny recording of accumulated drift for prior art capacitance manometers, has typically been accomplished by measuring such drift manually during routine zeroing. A judgment by the user then would then need to be made about the zeroing frequency, or whether the manometer has reached the end of its useful life.\nIn practice, users have typically not made such judgments due to various reasons including lack of time and/or complexity of factors involved. Instead, users have typically decided to zero or replace a manometer in reaction to an alarm on the related process tool or system or when they notice the manometer can no longer be zeroed (out of adjustment). In either case, the replacement or zeroing of the manometer causes an unexpected equipment failure problem (down time of the tool or system). Such unscheduled maintenance or replacements are something that must be dealt with outside of routine maintenance schedules, which is costly in terms of efficient use of the overall tool or system.\nThus, users have had to make a priori judgments about the frequency of zeroing. Determination of whether a manometer needs routine maintenance such as recalibration has generally not been done. Any recording of the manometer cumulative zero changes are typically also performed manually, an inconvenient process. In fact, the user generally only knows when a problem exists when there is a tool or system alarm.\nWhat is desirable, therefore, are systems, methods, and apparatus that address the limitations noted for the prior art by predicting and/or indicating when a capacitance manometer, such as those used for high-precision fluid delivery systems, requires zeroing, maintenance, and/or replacement, and automatically calibrating the system to correct for drift when possible."}
{"text": "Field of the Invention\nThe present invention relates to technologies for grasping with a robot hand a workpiece the three-dimensional position and orientation of which are recognized from among bulk workpieces, based on information about measurement by a sensor.\nDescription of the Related Art\nIn recent years, bulk picking technologies have been developed in which an individual object is identified from among a pile of workpieces by use of a vision system and then the three-dimensional position and orientation of the identified object are recognized to grasp the object with a hand attached to a robot in production lines in factories, etc. To realize continuous grasp operations of the robot in a short period of time in the foregoing technologies, it is important to preferentially select a workpiece that is graspable in a surface layer of the bulk where interference with a neighboring object is less likely to occur.\nJapanese Patent Application Laid-Open No. 2000-304509 discusses a method in which positional information about relative positions of a plurality of workpieces in the bulk is acquired and an overlapping state of the workpieces is acquired from the positional information to preferentially select a workpiece located in an upper area.\nIn the method discussed in Japanese Patent Application Laid-Open No. 2000-304509, the workpieces are rearranged in descending order according to their positions and orientations, and a workpiece located in a high position is preferentially selected. However, in a case where a region to be used as a reference for the position and orientation of a workpiece is different from a region that is to be grasped with a hand, or in a case where a workpiece with a longitudinal shape is graspable by grasping any of a plurality of regions that are significantly different in a longitudinal direction, it is not possible to preferentially select a workpiece that is graspable in a surface layer of the bulk. For example, FIG. 5 illustrates an example in which workpieces having a long shape in a particular direction are stored in bulk (five workpieces in this example) in a pallet. In FIG. 5, detected workpieces are specified by solid lines and referred to as workpieces A and B. Further, each of the workpieces is graspable by grasping either one of the regions referred to as target regions 1 and 2. In this case, the workpiece A can be grasped without interference with other objects in the surface layer of the bulk if the workpiece A is grasped by grasping the target region 1. However, in a case where a height based on the position and orientation of a workpiece is used as a reference, the workpiece B, which is more likely to interfere with a neighboring object, is selected. This decreases processing efficiency."}
{"text": "There are a number of approaches to solving the problem of modulating radio frequency (RF) carrier phase and frequency. The translational loop (TL), or offset phase-locked loop (PLL) uses a mixer in the feedback path of a transmit PLL to offset the output carrier frequency and an analog baseband I/Q data modulator in either the reference or feedback path of the transmit PLL to apply modulation. The primary drawback to the TL is that it requires well matched IF circuits that can occupy much more area than a design comprised primarily of digital blocks.\nAnother approach to modulating the output RF carrier is to apply the modulation via the feedback divider in a fractional-N PLL synthesizer. In this case a sigma-delta programmed multi-modulus divider is typically used to vary the loop output frequency and/or phase. The primary drawback to this architecture is that it requires careful calibration of the loop components and pre-emphasis filtering of the digital data in order to mitigate the distortion caused by the filtering in the PLL when the data bandwidth is greater than or comparable to the loop bandwidth. In some cases, this calibration can be more difficult by the fact that the fractional-N PLL bandwidth must be made narrower than the signal bandwidth in order to reject the shaped sigma-delta quantization noise and reference source noise.\nIn an integrated circuit with an integrated power amplifier (PA) in the modulation path, which causes temperature drift of loop components during a burst, the design requires either an adaptive pre-emphasis filter or a technique for adaptively adjusting the loop gain, both of which are non-optimal."}
{"text": "As examples of technologies concerning hip joint structures for legged mobile robots are known from Japanese Patent Publication No. 2592340 ('340) and Japanese Laid-Open Application Nos. 2001-62761 ('761) and 2001-150371 ('371). In '340 (especially at pages 4 and 5 and FIG. 2), there is taught a configuration in which the motors for driving the hip joints are disposed on the body side so as to reduce weight toward the distal ends of the legs and lower the inertial moment produced in the legs.\nThe teaching of '761 (especially at paragraphs 0053 to 0055 and FIG. 12) relates to a configuration in which the hip joints are provided with a parallel link mechanism with which the right and left legs are connected and the parallel link mechanism is operated at free leg footfall to move the leg upward so as to mitigate the ground impact force.\nThe teaching of '371 (especially at 0070 to 0086 and FIGS. 5 and 7) relates to a configuration in which each rotary shaft providing the degree of freedom to rotate about the yaw axis among the degrees of freedom of the hip joints is offset relative to the roll axis direction to avoid interference between the feet when the robot changes directions of movement.\nAside from the above, in the case where the body of the legged mobile robot is to be bent (i.e., bent forward or backwards), situations sometimes arise in which the desired amount of bending cannot be realized solely within the movable range of the respective shafts that rotate about the pitch axis and the roll axis of the hip joints.\nIt has therefore been proposed to increase the amount of bending of the body by dividing the body into an upper section and a lower section, connecting the two sections through a joint having a degree of freedom to rotate about the pitch axis and rotating the upper and lower sections relative to each other so as to realize an amount of body bending that is larger than that realized by the movable range of the hip joints (this can be seen by, for example, Kawada Industries, Inc., “‘Rise/Lie Actions’ Achieved by Humanoid Worker Robot,” [online], Sep. 19, 2002, Kawada Industries Homepage, topics, [retrieved May 2, 2003], Internet <URL: http://www.kawada.co.jp/general/topics/020919_hrp-2p.html>).\nHowever, the structure described on the aforesaid website has a problem in that the division of the body into upper and lower sections degrades the ability of the body to accommodate equipment internally.\nMoreover, the fact that the movable range of hip joints is deficient means not only that the amount of body bending is deficient but also simultaneously that the movable range of the legs cannot be increased. In the technology described on the website, the amount of bending of the body is increased by equipping the body with the joint, which does nothing to increase the moveable range of legs, so that no improvement is achieved in the degree of posture and gait freedom of the whole robot including its lower part."}
{"text": "Dental substances are often provided in devices allowing the substance to be dispensed directly to a desired location, for example on a dental pad or in a patient's mouth. Such dispensing devices typically have a chamber for holding the dental substance, an outlet, and a piston for extruding the substance from the chamber through the outlet.\nA variety of dispensing devices are designed for dispensation of relatively high viscosity dental substances, like for example dental filling materials. Some of those dispensing devices are configured for use with an applicator providing an extrusion force that is sufficient for dispensation of high viscosity dental substances. In dentistry several types of manually operated applicators are available which provide leverage for increasing manual forces to provide sufficient extrusion forces. Many applicators are designed as a reusable tool which forms part of a dental practice's basic equipment.\nA particular dispensing device often used to store, prepare and dispense a hardenable dental material mixed from a powder and a liquid is the so-called dental capsule in which the mixing of the powder and the liquid is performed within the capsule by shaking the capsule in a shaker.\nU.S. Pat. No. 8,893,925 discloses container for mixing and dispensing material. The container comprises a body having a main chamber, a dispensing nozzle, a liquid receptacle and a plunger. The liquid receptacle has a front portion arranged to break away upon pressure being applied by the plunger so that the plunger can traverse the entire length of the body. This enables a charge of material in the main chamber to be entirely dispensed through a frangible wall into the nozzle. The container is particularly envisaged for use in mixing and dispensing of dental materials.\nU.S. Pat. No. 8,968,000 discloses a mixing and application capsule for producing and discharging a dental preparation. To prevent a liquid loss, a mixing and application capsule for producing a dental preparation is proposed. The capsule has a body with a mixing chamber for receiving a mixing component and for mixing the dental preparation from the mixing component and a fluid. The capsule further has an outlet opening for discharging the dental preparation. A first plunger body which can be displaced in the body delimits the mixing chamber in the body. The first plunger body has a channel to guide the fluid from a cavity into the mixing chamber and a projection. The capsule has further a second plunger body which can be displaced in the body relative to the first plunger body. The second plunger body has the cavity to receive the fluid. The cavity is configured to receive the projection of the first plunger body.\nSuch capsules have a variety of advantages as, for example, a relatively long shelf life. However there is a still desire for a capsule with such a long shelf life, which is easy to use and which is nevertheless relatively inexpensive."}
{"text": "Water heating appliances and space heating appliances have undergone significant changes over the last several years. In the past, fuel was relatively inexpensive and water heaters and space heaters were designed for maximum reliability, simplicity and long life with only a small emphasis placed on fuel efficiency. This set of priorities has changed with the rising price of fuel. Over the last several years, water heaters and space heating furnaces have been designed with a much greater emphasis placed on fuel efficiency. These heaters are generally more complex than their less efficient ancestors.\nAn additional development has been the combining of water heating and space heating functions in a single appliance. This combination was developed to improve overall efficiency of the total water heating and space heating function within a given home, commercial building or vehicle. Moreover, savings on the space occupied by these appliances and installation expenses have been achieved by the combination of the two essential appliances into one high efficiency unit.\nOne such high efficiency combined water and space heating appliance is described in U.S. Pat. No. 4,541,410 to Jatana, the disclosure of which is incorporated herein by reference. Jatana describes a heating apparatus in which air is mixed with fuel and introduced into a blower which moves the mixture under pressure into the burner of a closed combustion chamber. The combustion chamber is contained within a tank containing water. The products of combustion exit the combustion chamber and pass through a helical tube of several turns within the body of water. The heat of combustion is extracted from the products of combustion by conduction through the walls of the combustion chamber and the helical exhaust tube. A high efficiency water heater results. The heated water from the water heater is also used to heat the air of a home or building by piping the hot water to a heat exchanger contained within the ducts of the home ventilation system. While this apparatus provides a highly efficient water and air heater, certain problems remain.\nIn order to maximize efficiency, the Jatana heater draws combustion air from outside of the home or building being heated. This requires piping or ducts from outside the building to the heater. Such piping or ducts are unique to each building in which the appliance is installed. Air to be mixed for combustion must be drawn through these ducts, and, due to variations in the length of the ducts, turns within the ducts and the like, the resistance of each individual installation to the flow of air will be different. This often requires that valves governing the air supplied to the furnace and gas supplied to the furnace be professionally and individually adjusted for each installation. Any changes in the resistance of the ducts to the flow of air caused by the introduction of foreign objects into the ducts, denting of the ducts or the like after adjustment may cause a change in the ratio of air to fuel in the heater, negatively effecting efficiency.\nThe burner in the Jatana device is a cylindrically shaped screen contained within a cylindrical combustion chamber. It has been found that the introduction of the air and fuel mixture into this burner under pressure sometimes results in a swirling circumferential motion leading to noisy operation.\nIn the Jatana unit and in many other high efficiency units several additional problems arise. Thus, air in excess of what is required for proper combustion of the fuel is often mixed with the fuel to be certain that the burner is never supplied with too little air, resulting in incomplete combustion.\nAlso, the burner of a high efficiency heater normally operates most efficiently at a given gas pressure and a specific energy input level. Variations in gas pressure and variations in energy input level cause reductions in efficiency. Therefore, the capacity, that is the input energy level, of a given furnace or heater cannot be easily adjusted while maintaining high efficiency.\nProblems of stability sometimes arise. Conditions in the combustion chamber change because of high outside wind, changes in gas pressure and the like and result in \"fast cycling\" and other problems which negatively effect efficiency. \"Fast cycling\" is the rapid and repeated changing of the furnace from the on state to the off state. It is often noisy and inefficient.\nThe present invention contemplates a new and improved apparatus and method which overcomes all of the above referred to problems and others and provides a water and space heating appliance of high efficiency, reliability, stability and quality."}
{"text": "The present invention relates generally to bath tubs and more particularly to a bathing tub for persons with some physical or mental impairment which generally results in the need for a bathing assistant. My invention is particularly useful for persons with impaired ambulatory ability, obesity, inadequate or unreliable judgment, persons who are susceptible to a loss of consciousness during the bathing process, and for those who have the desire for a safer, easier, more convenient method of bathing.\nIt will be appreciated by those skilled in the art that conventional bathtubs have a number of deficiencies in terms of their ability to serve the needs for physically or mentally impaired persons. First, it can be difficult for the physically or mentally impaired to get in or out of a conventional bathtub. Conventional bathtubs require a person to step over the side wall of the tub, an obstacle that may cause the physically impaired to trip and injure themselves in or about the tub. Also, the bathtub itself may have a slippery surface which is more difficult to navigate for the physically impaired, and non-slip surfaces and grab bars are often an inadequate solution for those who have minimal strength.\nShowers do not necessarily solve the problem because of the danger of slip and fall accidents and the inability of sick and infirm persons to stand for an extended period of time during showers. Also, showers do not lend themselves to assisted bathing since the assistant would get sprayed during the showering process.\nThere have been several attempts to address the problems associated with assisted bathing. One of the primary efforts has been the hoist-type systems which involves placing the person in a hammock-like structure, hoisting him over a tub of water and lowering him into the tub for the bathing process. The primary problems with the hoist-type systems are the fact that such devices are awkward and dangerous to manipulate, as well as being dehumanizing and humiliating to the patient.\nIn addition to the hoist-type systems, there have been developed a number of side access bathing tubs in which the patient lowers or moves the side of the tub out of the way, sits into the tub, raises or moves the side of the tub back in place, has the tub filled with water and then engages in the bathing process. Problems with devices of this nature involve side door maintenance, and sealing to avoid leaks in the tub, delays in the tub filing process, having to empty the tub if the person has to get out of the tub before the bath is completed, etc.\nAnother approach has been to provide side access with a tub chamber that tilts for access. However, these devices as well as all other prior art side access bathing tubs cannot be pre-filled and the bather must suffer the discomfort while the tub is filling of sitting in a potentially cold tub during the extended time that it takes to fill the tub. Further, the tub surfaces cannot be pre-warmed by having the water in the tub immediately prior to the bather entering the tub and the bather can be scalded or chilled by the water as it is filling the tub or irritated or injured by undiluted additives added to the bath water while the bather is in the tub.\nPractically all of these prior art efforts at solving problems associated with assisted bathing for the mentally or physically impaired involve structures that cannot be quickly drained and re-filled with the same bathing liquid resulting in a waste of water when the bather has to get out of the tub in the middle of the bathing process, for example to go to the toilet, involve dangers in egress and ingress in the tub, or are inconvenient to use. Further, these prior art devices do not deal with the problems associated with unattended immobility of the bather, loss of consciousness and the like that may result in drowning."}
{"text": "1. Field of the Invention\nThe present invention relates to digital broadcasting, and more particularly, to a digital television signal and a method and apparatus of processing a digital television signal.\n2. Discussion of the Related Art\nGenerally, a program and system information protocol (hereinafter, PSIP) is a protocol for channel tuning and broadcast schedule transmission in ATSC (advanced television systems committee) as the digital broadcasting standard in terrestrial and cable digital broadcasting environments. In addition, the PSIP is a standard protocol for a transfer of tables included within packets transferred by a multiplexed transport stream.\nThe PSIP defines many PSIP tables for the various purpose, and each PSIP table includes one or more sections which are similar to those of the program and system information (PSI) tables defined by MPEG (moving picture experts group). Each section included in a PSIP table includes a header containing basic information of the section and a body containing the actual data constituting the section. A broadcast receiver uses the basic information included in headers of various sections received from a broadcast transmitter in order to select (or detect) one or more pertinent sections, and this is often referred to as section-filtering. An example of the basic information is a table identifier which identifies a PSIP table.\nHowever, a certain group of PSIP tables (e.g., a group of event information tables or extended text tables) have a common table identification. Therefore, the digital television (DTV) receiver is not able to distinguish them simply by section-filtering. In this case, the DTV receiver must parse bodies of the PSIP sections in order to detect a pertinent section, and this greatly increases the signal processing time and decreases the system's overall efficiency."}
{"text": "1. Field of the Invention\nThis invention relates to a ski sled which combines the thrill of skiing with the safety and ease of operation of a sled wherein a seated operator can control the direction of the ski sled by tilting the skis and can operate a brake to slow down or stop the ski sled.\n2. Description of the Prior Art\nThose interested in winter sports often start sledding as children. A sled is easy and safe to operate but a typical sled with steel runners is hard to steer and does not function very well on many surfaces such as deep, soft snow. Therefore, most people who want to continue with outdoor activities learn to ski because that offers greater variety and more excitement. Not everyone can learn to ski, however, particularly where physical handicaps are involved. Further, some persons who may not have time to learn to ski would still like to periodically enjoy riding down a hill like a skier but without having to stand up and balance on a pair of skis.\nAttempts have been made in the prior art to provide such an apparatus. A ski sled is shown in the Nuss U.S. Pat. No. 3,107,923 issued Oct. 22, 1963. That ski sled has a mechanism for twisting the front portions of the skis to steer the ski sled. No braking mechanism is shown. A snow shuttle is disclosed in the Chilzer U.S. Pat. No. 4,349,209 issued Sep. 14, 1982. That snow shuttle has a single, rather wide runner at the rear and a front runner which can be pivoted with respect to the rear runner for steering purposes, the rear runner having a braking member associated therewith. The Chilzer snow shuttle is more like a conventional sled and would not offer the thrills skiing. A ski maneuvering apparatus having brackets connecting the skis for steering the skis, but without a seating platform or braking mechanism, is shown in the Zepkowski U.S. Pat. No. 4,357,036 issued Nov. 2, 1982. Various other mechanisms can be found in the prior art, but all have drawbacks similar to those mentioned above."}
{"text": "1. Field of the Invention\nThis invention relates to a recycle developing process used in an electrophotographic apparatus such as a copying machine or a printer.\n2. Description of the Prior Art\nGenerally, image formation in electrophotography is carried out by charging (main charging) the surface of a photosensitive material, imagewise exposing the charged material to form an electrostatic image on the surface of the photosensitive material, developing the image with a developer filled in a developing vessel to form a visualized toner image, transferring the toner image to a predetermined paper, and removing the toner remaining on the photosensitive material after the transferring by using a device such as a cleaning blade to complete one cycle of image forming step.\nA typical developer is, for example, a two-component magnetic developer composed of an electroscopic toner made of a colored resin composition and a magnetic carrier. The developer is delivered to a developing zone in the form of a magnetic brush by means of a developer conveying sleeve provided in a developing vessel, sliding the magnetic brush with the electrostatic image on the photosensitive material, and adhering the toner to the electrostatic image to thereby perform development.\nRecently, many recycle developing processes in which a toner removed and recovered by cleaning is again recirculated into a developing vessel and again used for development with an object of reutilizing toners have been proposed, and have been applied to actual electrophotographic apparatuses. This recycle developing process is applied to inexpensive low-speed machines using organic photosensitive materials (OPC) in general.\nIn this recycle developing process, when a toner in the starting developer filled in the developing vessel is consumed and reaches a concentration of a predetermined level or below, a virgin toner is replenished from a feeding hopper, and a toner recovered by cleaning is also replenished.\nHowever, in the above-mentioned recycle developing method, there is a problem in that the properties of a recovered toner to be used again after recovering by cleaning are different from the properties of a toner contained in the starting developer or those of a virgin toner supplied to a developing vessel.\nFor example, the surface of the toner is surface-treated with a treating agent such as silica or alumina so that its properties such as flowability may be held stably. However, the toner which is supplied for development is adhered to the surface of the photosensitive material and thereafter is recovered by cleaning, the surface treating agent is removed or embedded in the toner particles due to an external force of cleaning or to a force exerted after recovering in the step of conveying into the developing vessel. Thus, the toner has an extremely decreased flowability. Accordingly, as the recovered toner is replenished into the developing vessel, the flowability of the developer decreases. Especially when some degree of development is conducted, the toner in the developer contained in the developing vessel all becomes a recovered toner, and the properties of the toner become very much different from the initial developing agent, and it becomes difficult to perform development stably.\nIn the developing vessel, a toner concentration sensor is provided so that the toner concentration (T/D) of a developer composed of a toner and a carrier may be controlled within a predetermined range. Controlling of this toner concentration is carried out by utilizing the variation of the toner concentration in the developer corresponding to the permeability of the developer. The permeability of the developer is detected by the toner concentration sensor and the toner is replenished into the developing vessel according to the output value of the sensor.\nChanges in the properties of the developer by the replenishing of the recovered toner adversely affect the controlling of the toner concentration. For example, curve A in FIG. 3 shows the relation between the output (corresponding to the permeability of the developer) of the concentration sensor in the starting developer and the toner concentration (T/D). According to this curve, if the threshold value of ON-OFF of toner replenishing is set at a sensor output value 3V, the toner is replenished to the developing vessel when the toner concentration becomes 3.5% or below. However, when the recovered toner is replenished into the developing vessel, the properties vary and the relation between the output of the concentration sensor and the toner concentration changes to, for example, curve B. Therefore, by the above-mentioned setting of the threshold value, it will become difficult to maintain the toner concentration at a predetermined level."}
{"text": "There are a number of U.S. patents that disclose electronic apparatus for printing indicia on labels, some of these are restricted to hand held units and others that disclose tabletop units. Hand held label printers, such as disclosed in U.S. Pat. No. 6,113,293, and tabletop printers, such as disclosed in U.S. Pat. Nos. 6,266,075 and 5,078,523, include the same general combination of elements, a print head, means for feeding label media to be printed past the print head, a microprocessor, a read only memory programmed with appropriate instructions to operate the microprocessor, a random access memory, a keyboard with letter, number, and function keys for the entry of alphanumeric information and instructions concerning the indicia to be printed, and a visual display such as a light emitting diode (LED) or liquid crystal display (LCD) unit to assist the operator in using the printer. In a hand held printer, these components may all be enclosed in a single housing.\nThe label media comprises a series of labels that are attached to a carrier strip. The carrier strip is fed through the printer and legends, alphanumeric characters, and other indicia, are printed on the labels. The labels are then removed from the carrier and attached to the objects needing identification. As there are many types of label applications, there are many combinations of labels and carrier strips that provide labels of varying sizes, colors and formats.\nA particular type of print head employs thermal transfer printing technology. Thermal transfer printing uses a heat generating print head to transfer a pigment, such as wax, carbon black, or the like, from a thermal transfer ribbon to a label media. By using digital technology, characters are formed by energizing a sequence of pixels on the print head which in turn melts the wax or other pigment on the ink ribbon transferring the image to the label media.\nIn a known thermal transfer printer such as a label printer, label media and ink ribbon are simultaneously fed past the print head by a platen roller in an overlay relationship between the print head and the platen roller. The platen roller is rotatably driven by a drive mechanism that may also rotatably drive ink ribbon take up and supply spools to maintain tension in the ink ribbon.\nIn a cartridge-based printing system, such as disclosed in U.S. Pat. No. 6,113,293, it is desirable to have a consistent label media path. In order to accomplish this, many cartridge-based printing systems have the label media path defined by a point tangent to the outside diameter of a roll of label media, such as a label media drive roller. This method, however, presents a problem as the label media is consumed. In particular, as the label media is consumed the diameter of the roll decreases and the beginning point of the label media path changes.\nThis problem is solved by pivotally mounting the label media, such that as the diameter of the label media decreases, the label media pivots to maintain the label media in contact with the label media drive roller. Unfortunately, as the label media is jostled or jarred, the label media can lose contact with the label media drive roller and retract into the cartridge. As a result, the end of the label media may be difficult to retrieve from the cartridge for use. Accordingly, a need exists for an improved printer cartridge that can carry label media that remains in contact with a label media drive roller even when jarred or jostled."}
{"text": "Multi-media integrated circuits (ICs) generally need to transpose a large amount of 2D data. High definition (HD) JPEG pictures taken from a camera need to be rotated before being displayed by a DVD recorder/player. Optical discs, such as Blue-Ray or HD-DVD, have 2D error correction for a data stream. Data is stored in memory in an order defined by the bitstream. However, the data needs to be accessed in column order.\nConventional approaches tend to either use embedded processors to move the data around or to implement complicated dedicated hardware to rotate the data in a matrix. The software approach implemented in an embedded processor is slow. The amount of dedicated hardware tends to grow exponentially with the size of the data in the matrix.\nThere are several applications which need to transpose a 2D data matrix at a high speed. One such application includes image rotation. With image rotation, millions of pixels of a still image are captured by a camera. If an image is taken with the camera rotated 90 degrees, then the image needs to be rotated before being shown on a display monitor. The rotation of an image is normally done by a DVD player/recorder. Such a rotation has to be completed within a reasonable time, often less than 1 second. As the resolution of cameras increases, additional dedicated hardware is needed to meet predetermined performance goals.\nAnother conventional approach involves optical data error correction. Such an approach receives data from the optical disc, such as blue ray or an HD DVD system. The received data needs error correction (ECC) to be performed. Part of the ECC process involves mathematical operations based on data columns, as opposed to an operation which is based on data stream order. Graphic data manipulation includes graphic data which needs a rotation operation.\nIn previous generations of DVD systems, the rotation operation was performed by an embedded processor where data is rotated by software. In next generation DVD systems, a hardware rotation engine is needed to support a high speed rotation for a large amount of data on the fly. However, rotation data needs to read data from the same column position. In older generation chips, multi-port random access memories (RAMs) or registers with hardwired multiplexers were used. The use of multi-port RAMs or registers with hardwired multiplexers is not practical as the size of a data matrix continues to increase, often exponentially.\nIt would be desirable to provide a method and/or apparatus for an efficient and/or high speed two dimensional data transpose engine for a SOC in a manner that may be implemented in a low cost and/or reduced size integrated circuit."}
{"text": "For instance, devices for lifting and lowering a vehicle seat seating surface are roughly classified into stepless lifter devices and stepped lifter devices which lift the seating surface steplessly and stepwisely, respectively."}
{"text": "Digital recording equipment enables users to capture and store digital media content items, for example, video, images, audio, and the like. Software applications enable users to perform operations on the captured and stored digital media content. Such operations include viewing, editing, and the like. For video segments, a software application's editing operations can include deleting portions of the video segment, combining two or more video segments, overlaying content on a stored video segment, and the like. Some software applications allow users to associate audio with video content. After such associating, as the video content is played, the associated audio content is also played, e.g., as background music to the video content. In some scenarios, multiple audio tracks can be associated with video content. For example, a first audio track can include audio that was captured along with the video content and a second audio track can include background music that is associated with the video content after the video content is captured."}
{"text": "Crude oil produced from geological formations can contain various amounts of water. Water and crude oil are naturally non-miscible; however, when naturally occurring interfacial active compounds are present, these compounds can aggregate on the water and oil interface and cause water to form droplets within the bulk oil phase. During crude oil lifting through production tubings, the water and oil encounters an increased mixing energy from rapid flow through chokes and bends. This additional mixing energy can emulsify the water and oil. Such an oil external, water internal two-phase system is commonly referred to as a crude oil emulsion, which can be quite stable. The presence of water in crude oil, however, can interfere with refining operations, induce corrosion, increase heat capacity, and result in reduced handling capacity of pipelines and refining equipment. Therefore, the crude oil that is to be shipped out of the oilfield should be practically free of water and usually has a maximum water content limit of about 0.5 to 3% by total weight, depending on the type of crude and oil company.\nThe emulsified water can also contain various amounts of salts. These salts are detrimental to crude oil refining processes due to potential corrosion in the refinery. In crude oil refining, desalting techniques comprise the deliberate mixing of the incoming crude oil with a fresh “wash water” to extract the water soluble salts and hydrophilic solids from the crude oil. Primary dehydration of the crude oil occurs in oil field water oil separation systems such as “free water knock out” and “phase separators.” Quite often, these systems are not adequate for efficient separation due to factors such as over production, unexpected production changes, and system underdesigns. In these cases, emulsion-breaking chemicals are added to the production processes to assist and promote rapid water oil separations.\nCommonly used emulsion-breaking chemicals or demulsifiers include alkylphenol formaldehyde resin alkoxylates (AFRA), polyalkylene glycols (PAG), organic sulfonates, and the like. These active ingredients are typically viscous and require a suitable organic solvent to reduce the viscosity of the demulsifier blend. Accordingly, there is an ongoing need for new, economical, environmentally-friendly, and effective chemicals and processes for resolving emulsions into the component parts of water and oil or brine.\nFor these reasons it is desired to have a demulsifier that does not require an organic solvent. It is also desired to have a demulsifier composition capable of resolving water external, water internal, and complex emulsions while having a broad dosage range."}
{"text": "1. Technical Field\nThe present invention relates generally to fiber optic cables, and in particular relates to fiber optic drop cable assemblies for deployment on building walls.\n2. Technical Background\nFiber optic communications has experienced explosive growth. In just a few years, the transmission of communication signals for voice, video, data, and the like has soared, and more growth is planned as fiber optic communication technology improves and networks expand to provide greater access.\nFiber optic cables are the backbone of fiber optic communication systems. Fiber optic cables carry optical fibers and other cable elements, which are protected from the external environment by an external jacketing.\nFiber optic cables that carry optical signals to a home or other locations from a connection point (e.g., a local splice box or “fiber distribution terminal” or FDT) on the distribution cable in so-called “fiber-to-the-X” (FTTX) networks are referred to in the art as “drop cables.” In office buildings and apartment buildings (referred to in the art as “multiple dwelling units” or “MDUs”), drop cables are preferably deployed internal to the building, usually through ductwork or molding. However, for certain buildings, and in particular older buildings, such deployment can be time consuming and expensive. Further, in many older buildings, there may actually be no practical way to internally route cabling through the infrastructure. One present alternative to internal cable deployment is to deploy copper cable on the outside of the building. However, this approach is expensive and does not provide the telecommunication performance offered by fiber optic cables.\nDeploying fiber-based drop cables on the outside wall of a building would be a preferred cabling solution in terms of telecommunications performance, material cost, and speed of deployment. However, such a solution is made very difficult by the bending tolerances of typical fiber-based drop cables, and is complicated by the need to ground the steel messengers used to offload mechanical loads and strains that would otherwise be borne by the fiber optic cable."}
{"text": "In the medical and veterinary fields bodily fluid samples such as urine are quite frequently obtained to assist in the diagnosis of ailments. Since this invention is particularly useful in the medical field it will be discussed in that context, however, it should be understood that it is applicable to the veterinary field. On occasion a fluid sample, such as a urine sample procured over a 24 hour period, could be used for several things. Initially, it may be important to determine the volume of fluid passed by the individual during a 24-hour period. Secondly, body fluid, such as urine obtained over a 24 hour period, may be helpful in determining the presence of an ailment such as diabetes and the like. It is also common that more than one test may be performed on the 24 hour fluid sample. While it would be possible to divide the sample after it is fully obtained from the patient, certain tests require different preservatives in order to maintain the fluid in condition for testing. If the preservatives in each of the tests were identical, then division of the sample at the end of the 24 hour period would be simple. Since preservatives quite frequently are mutually exclusive division of the sample after collection is precluded unless a third container is used for collection and immediate division.\nBecause of the mutual exclusivity of the preservatives used in 24 hour collection tests, it is common practice for physicians when obtaining urine, to require a patient to collect fluid not over a 24-hour period but over a 48-hour period. While this provides the physician with the necessary samples for testing, it works an inconvenience on the patient in that the patient's movements are restricted for twice the period of time.\nIt might be suggested that the urine be divided at the time the sample is passed. This could be a workable solution if the patient could always be relied upon to divide the sample obtained at each passing such that one half was placed in one collection bottle and the other half placed in a second collection bottle. In the hospital situation where the patient provides the sample to a nurse, the physician is dependent upon the nurse to divide the sample throughout the 24-hour period. While this seems to be a viable solution from the standpoint of the nursing profession, it is still subject to human error thus an entire sample may have to be discarded if a descrepancy occurs during the sampling.\nAccordingly, this invention provides a simple inexpensive fluid proportioner which may be used in association with 24 hour body fluid sample.\nIt is a further object of the invention to provide a fluid proportioner that may be utilized to divide samples into two relatively equal sample sizes over an elapsed period of time.\nIt is still a further object of this invention to provide a fluid proportioner that is inexpensive.\nIt is still a further object of this invention to provide a fluid proportioner that will withstand the hostile environment of a sterilization process.\nIt is also an object of this invention to provide a fluid proportioner that may be used in conjunction with standard sample containers and funnels.\nThe invention is a fluid proportioner that includes a first tubular member having an open end and a closed end. At least two second tubular members each having a cross sectional area less than the cross sectional area of the first tubular member are affixed to and open into the closed end of the first tubular member."}
{"text": "The present invention relates in general apparatuses and methods for facilitating thermal analysis of a data center, and more particularly, to apparatuses and methods for thermally simulating one or more respective electronics subsystems, electronics racks or rows of electronics racks to be disposed within a data center.\nThe power dissipation of integrated circuit chips, and the modules containing the chips, continues to increase in order to achieve increases in processor performance. This trend poses a cooling challenge at both the module and system level. Increased air flow rates are needed to effectively cool high power modules and to limit the temperature of the air that is exhausted into the computer center.\nIn many large server applications, processors along with their associated electronics (e.g., memory, disk drives, power supplies, etc.) are packaged in removable drawer configurations stacked within a rack or frame. In other cases, the electronics may be in fixed locations within the rack or frame. Typically, the components are cooled by air moving in parallel air flow paths, usually front-to-back, impelled by one or more air moving devices (e.g., fans or blowers). In some cases it may be possible to handle increased power dissipation within a single drawer by providing greater air flow, through the use of a more powerful air moving device or by increasing the rotational speed (i.e., RPMs) of an existing air moving device. However, this approach is becoming problematic at the rack level in the context of a computer installation (e.g., data center).\nThe sensible heat load carried by the air exiting the rack is stressing the ability of the room air conditioning to effectively handle the load. This is especially true for large installations with “server farms” or large banks of computer racks close together. In such installations not only will the room air conditioning be challenged, but the situation may also result in recirculation problems with some fraction of the “hot” air exiting one rack unit being drawn into the air inlet of the same rack or a nearby rack. This recirculation flow is often extremely complex in nature, and can lead to significantly higher rack inlet temperatures than expected. This increase in cooling air temperature may result in components exceeding their allowable operating temperature or in a reduction in long term reliability of the components.\nAddressing thermal imbalances within a data center is often an expensive and time consuming operation. Therefore, there is a need in the art for apparatuses and methods which facilitate thermal and energy based design, analysis and optimization of electronics equipment of a data center."}
{"text": "Children have enjoyed playing with toys, such as dolls, action figures, etc., including vehicles such as toy train sets for many years. Toy train sets come in many different forms, such as model railroad sets, remote controlled sets, and wooden sets.\nExisting toy train sets include tracks and support pieces that form configurations that remain on a single level. Some toy train sets can be built to form configurations on two levels, i.e., some tracks positioned over and above other tracks that are supported by a support surface. Any track layouts having configurations of more than two levels can be unstable, unuseable, and non-entertaining due to the track collapsing."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a socket connector, and more particularly to a Zero Insertion Force (ZIF) socket connector for electrically interconnecting a Central Processing Unit (CPU) with a Printed Circuit Board (PCB).\n2. Description of Related Arts\nU.S. Pat. No. 6,267,615 issued to Lin on Jul. 31, 2001 discloses an electrical contact comprising a main portion, a soldering portion, a pair of wing portions extending laterally and forwardly from two edges of the main portion and a pair of contact portions each formed at a free end of the wing portion. However, such kind of electrical contact consumes a large quantity of material in manufacture because the electrical contact takes a large width in a plane view before the wing portions are folded. Therefore, each two adjacent electrical contacts have a larger distance therebetween before the contacts are cut off from the belt. Usually, the adjacent electrical contacts usually have double pitch which is two times of a distance defined between the adjacent passageways of the insulative housing. In an assembling of the electrical contacts into the insulative housing, two or three times are usually needed to insert the same row of electrical contacts, which is rather troublesome and time-using.\nU.S. Pat. No. 6,471,534 issued to Lee on Oct. 29, 2002 discloses an electrical contact includes a retention portion, a soldering portion extending from the retention portion, a first arm and a second arm extending upwardly from the retention portion and beside each other. The first arm has a beam section, a curved section and a free end. The second arm includes a leg section and a curved contacting section. The free end and the curved contacting section cooperate with each other to press therebetween and electrically connect with a pin of an electronic package received on the cover. However, the first arm has only one curved section that it makes a cut-off section of the material contacting with the inserted pin of the CPU. The pin of the CPU is easily scraped by the cut-off section because of the cut-off section is rough. After several times of insertion and removal, the pin of the CPU may be damaged and unable to use.\nHence, an electrical connector with a material-saving electrical contact for preventing damage to the pin of the CPU are desired."}
{"text": "A metallic spring finger is a device which is capable of biasing an object or providing a holding force on the object in order to maintain the object at a fixed position relative to a main body. FIG. 1 shows a side view 20 of a metallic spring finger 22 attached to a main body 24. As shown, the metallic spring finger 22 includes (i) a V-shaped end portion 26, and (ii) a substantially flat middle portion 28 which connects the V-shaped end portion 26 to the main body 24. Here, the overall length (L) of the entire spring finger 22 from the main body 24 to the end of the spring finger 22 is based on the length contributions of both the middle portion 28 and the V-shaped end portion 26 (see FIG. 1).\nSuch a metallic spring finger 22 applies a reaction spring force (F) in response to a displacement (D), e.g., due to displacement by an object 29. The dashed lines show the metallic spring finger 22 in a new position after being displaced from its original position. In general, the spring force (F) provided by the metallic spring finger 22 is a function of a spring constant (k) and the displacement (D).F=k*D  (1).Equation (1) illustrates the function for determining the spring force (F) provided by the metallic spring finger 22 in FIG. 1.\nFIG. 2 is a graphical illustration as to how the spring force (F) of the above-described conventional metallic spring finger 22 increases linearly as the displacement (D) increases. The slope of the curve in FIG. 2 is equal to the spring constant (k).\nIt should be understood that due to the above-described property of the metallic spring finger 22, the metallic spring finger 22 is suitable for positioning or holding an object relative to the main body 24. For example, as shown by the solid lines in FIG. 1, the metallic spring finger 22 may initially extend from the main body 24 at a 90 degree angle. Then, when an object 29 is placed at a location 30 between the metallic spring finger 22 and the surface 32 upon which the main body 24 is mounted resulting in the metallic spring finger 22 being pushed in a direction away from the surface 32 to the new position which is slightly greater than 90 degrees (shown in phantom in FIG. 1), the metallic spring finger 22 responds by providing the force (F) against the object to hold the object at the location 30."}
{"text": "During the past several years, the telecommunication industry has enjoyed explosive growth. The industry has strained to meet the demand for increasing communication bandwidth. Global, national and regional telecommunication techniques primarily include: (1) telephone networks providing voice, data and FAX transmission using twisted pair wire, coaxial cable, fiber optics, microwave systems, and RF networks; and (2) television networks providing television through RF transmission, and coaxial cable, and (3) internet data utilizing some or all of the specified data transmission media. Television, telephone and data communication are also currently being provided through satellite-based systems. Growth in non-voice data communications including high-speed image and video data establishes the need for data communication at increasingly higher speeds than that required for voice communication alone.\nRadio communication permits the user to be mobile. It does not require expensive wiring connecting the communication equipment. The problem with radio communication is that the available radio bandwidth is limited. A solution to limited bandwidth is to establish separate geographical cells which allows reuse of available bandwidth for each separate cell. By establishing a large number of independent cells, the number of users that can share the common bandwidth is increased. Nevertheless, since the RF data transmission is essentially broadcast, the permissible density of common-band cells, and likewise the maximum density of users that can be accommodated, is necessarily limited.\nThe need for alternative technologies for last mile and last few miles connectivity solutions is on the increase. In the past telephone and cable systems have generally operated on a regulated monopoly basis. Currently, however, the federal, state and local governments in the United States are encouraging competition in the provision of these services. Still, local telephone and cable companies are reluctant or charge heavily to share their installed infrastructure, and the installation of new cable or fiber connections can be prohibitively expensive, disruptive or otherwise not possible. In many developing countries there is no significant wired communication infrastructure in place and installing a wired infrastructure would be expensive and disruptive. Certain events such as the Olympic Games and the Super Bowl create temporary need for greatly expanded communication in a region. Disasters such as major ice storms or hurricanes can disrupt existing communications creating a need for temporary high-bandwidth communication equipment until the existing system can be repaired.\nTechniques for providing free space optical communications are known. (See xe2x80x9cA Brief History of Free-Space Laser Communicationsxe2x80x9d by David L. Begley in Selected Papers on Free-Space Laser Communications, David L. Begley, ed., SPIE Optical Engineering Press, 1991.) For example free space laser communication has been proposed for satellite to satellite communication, where the laser beam can provide extremely high bandwidth over long distance point-to-point links. Satellite communications provide a nearly ideal venue for the application of laser communication technology, where the beam path is above the earth\"\"s atmosphere. By contrast, ground-based free space laser communication has generally been constrained to relatively short distances because of the adverse effects of atmospheric conditions such as turbulence.\nLaser sources are capable of providing extremely narrow beams, with footprints of tens of centimeters at kilometer distances easily achievable. These extremely narrow beams permit an almost unlimited number of local links to operate at the same carrier frequency without interfering with one another and yet requiring low transmission power levels to achieve good signal-to-noise. The difficulty with arbitrary reduction of transmitted beam size is in maintaining alignment between transmitters and receivers. Once the beam is reduced below a few hundred microradians, even the most minor vibration or other perturbation will be sufficient to cause the transmitted beam to entirely miss the remote receiver.\nWhat is needed is a laser communication system with transceivers which can be installed and easily aligned and will remained aligned.\nThe present invention provides a free space laser communication system which dynamically maintains link alignment. The system comprises two laser transceivers positioned a distance apart. Each transceiver includes the following components. A beacon laser is provided for transmitting a beacon beam at a first wavelength in a beam having a first divergence. A signal laser is provided for transmitting a signal beam carrying information to be transmitted to the other transceiver. A special telescope system is provided for collecting both the incoming beacon beam and the incoming signal beam from the other transceiver and for transmitting an outgoing signal beam. A tracker is provided to maintain dynamic alignment. A high-speed detector is provided to receive the data transmitted from the remote terminal.\nThe optical system images the pupil of the telescope onto a tracking mirror, which pivots in tip and tilt. The received beacon beam is reflected from the tracking mirror and focused on a CCD array that monitors the location of the remote beacon source and provides feedback to the tracking mirror. Based on this feedback, the tracking mirror tips and tilts as necessary to maintain the received beacon beam on a predetermined position of the CCD array, while simultaneously directing the received signal beam through a small aperture to the high-speed signal detector. The small aperture isolates the incoming signal beam from spatially distinct background radiation. The outgoing laser beam from the signal laser is also reflected from the tracking mirror and out of the telescope system toward the remote transceiver. In a preferred embodiment the signal laser for both transceivers is a single mode diode laser operating at 850 nm. The divergence of the signal beam from each transceiver is approximately 200 microradians, providing a footprint at a two-mile transceiver separation of approximately 2 feet. The beacon laser is a diode laser operating 785 nm. The beacon beams are given much larger divergences on the order of one degree to provide for tracking over a wide dynamic range. Preferred embodiments permit communication at high data rates of up to 2.5 gigabits per second at ranges in excess of six miles in clear weather. The dynamic and automatic tracking features of the present invention permit the system to be installed and made operational in just a few minutes. In case alignment is disrupted for any reason the system automatically realigns itself as soon as the cause is removed.\nThe present invention provides a laser-based telecommunication system, which can be installed and aligned easily, efficiently and which will remain aligned for extended periods of operation. Additionally, this system can be operated over distances ranging from hundreds of meters to greater than ten kilometers, with communication speeds in excess of gigabits per second and with high link reliability in the face of local perturbations and path atmospheric turbulence."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to the design of collision detection and warning systems. More specifically, embodiments of the present invention relate to a method and system for facilitating collision warning.\n2. Related Art\nRecent developments in sensing and mobile computing technologies have allowed consumers to benefit from computing in day-to-day applications, which previously has only been available in expensive, stationary systems. One such application is vehicle collision warning. Driver safety can be significantly improved with present-day sensing technologies combined with cost-effective computing power, mobile communication capabilities, and ubiquitous positioning systems. Equipped with these technologies, vehicles can access much better information about the road ahead, and can avoid many more accidents.\nTo accurately predict collisions and provide effective warning, a collision warning system is often required to process a huge amount of information from different sources. This information overload is further exacerbated as sensors grow more sophisticated and more information about a vehicle's surroundings becomes available. For example, such information may come partly from better sensors in vehicles, partly from intelligent intersections which can track the movements of many principals (e.g., vehicles, pedestrians, and bicyclists), and partly from other principals which carry intelligent devices for tracking and communicating their movements to peer vehicles.\nHence, it is a significant challenge to process this large flow of information in real time and to give accurate and useful warnings of a potential collision."}
{"text": "In order to prevent a safety accident during driving, a forward collision warning system calculates a safety distance from a forward vehicle, and generates an alarm to warn a driver when a collision is likely to occur. The forward collision warning system senses the forward vehicle through a sensor or camera, and calculates a relative distance or relative velocity between the ego vehicle and the forward vehicle. Then, the forward collision warning system calculates a time to collision (TTC) from the relative distance and the relative velocity, and warns of a collision with the forward vehicle using the TTC.\nFor example, as illustrated in FIG. 1, an image taken by a camera mounted on a vehicle traveling on the road 10 shows a forward situation. When a forward vehicle 12 is present in the image, the forward vehicle 12 may be detected through the following process. First, the forward collision warning system acquires a forward image through the camera, and performs inverse filtering on the acquired image in order to distinguish the forward vehicle 12 from other objects. Then, the forward collision warning system recognizes the shape of the forward vehicle 12, and determines whether the forward vehicle 12 is a vehicle. When it is determined that the forward vehicle 12 is a vehicle, the forward collision warning system recognizes the bottom of a rear wheel or the end of a lower shadow area of the forward vehicle 12 as a distance reference point, and calculates a relative distance to the forward vehicle 12 based on the distance reference point.\nIn FIG. 2, the left drawing includes a coordinate axis indicating an actual position of the forward vehicle 12 relative to the ego vehicle 14, and the right drawing includes a coordinate axis in an image taken by the camera of the ego vehicle 14. Here, the origin point corresponds to a vanishing point. Since the distance reference point is set to the bottom of the rear wheel or the end of the lower shadow area of the forward vehicle 12 as described above, ‘Y=0’ in the coordinate axis of the forward vehicle 12. Thus, the forward collision warning system can simply calculate the relative distance to the forward vehicle as expressed by Equation 1 below.\n                    Z        =                  λ          ⁢                      h            y                                              [                  Equation          ⁢                                          ⁢          1                ]            \nIn Equation 1, Z represents the relative distance between the ego vehicle and the forward vehicle, λ represents a focal distance of the camera, h represents the height of the camera, and y represents the height of the distance reference point on the camera.\nWhen the relative distance calculation is repeated in successive images, the relative velocity may be calculated through Equation 2 below.\n                    v        =                                            Z                              t                +                                  Δ                  ⁢                                                                          ⁢                  t                                                      -                          Z              t                                            Δ            ⁢                                                  ⁢            t                                              [                  Equation          ⁢                                          ⁢          2                ]            \nIn Equation 2, v represents the relative velocity of the ego vehicle relative to the forward vehicle, Zt represents a relative distance at time t, and Zt+Δt represents a variation of the relative distance depending on a time variation Δt.\nWhen a TTC obtained by dividing the relative distance by the relative velocity is less than a threshold value as expressed by Equation 3 below, the forward collision warning system warns of a collision risk.\n                    TTC        =                              Z            v                    <                      TTC            threshold                                              [                  Equation          ⁢                                          ⁢          3                ]            \nIn Equation 3, TTC represents a time to collision, and TTCthreshold represents the threshold value of the TTC.\nThe conventional forward collision warning system can easily acquire the relative distance. As illustrated in FIG. 3, however, when the conventional forward collision warning system cannot detect the rear wheel of the forward vehicle because the forward vehicle is too close to the ego vehicle, the conventional forward collision warning system may not acquire the distance reference point or cause a misrecognition. Such a situation may frequently occur when the vehicle is driven at low speed. Furthermore, since a collision warning is not generated in time, a collision accident may occur.\nKorean Patent Publication No. 10-2013-0007243 discloses a forward collision warning method and system capable of successfully recognizing a forward vehicle under various illumination intensities or road conditions. Referring to the related art document, the forward collision warning method and system minimizes the load of a CPU through vehicle recognition and tracking. However, a vehicle tracking unit of the related art document detects the bottom of a forward vehicle recognized by a forward vehicle recognition unit and sets the bottom of the forward vehicle to a feature point of the forward vehicle. Therefore, when the forward vehicle is too close to the ego vehicle as illustrated in FIG. 2, the forward collision warning method and system cannot acquire a feature point, and may fail to generate a collision warning."}
{"text": "The popularity of multiple wireless communication technologies for handheld platforms has created a need to integrate wireless communication technologies on a single wireless communication device. However, the frequency bands of some of these technologies are close enough to result in interference. For example, the un-licensed 2.4 GHz Industrial, Scientific and Medical (ISM) frequency band is adjacent to some of the bands used by Mobile Wireless Standards (MWS) technologies to result in adjacent channel interference. In many electronic devices such as smartphones, both ISM and MWS technologies are implemented in the same device. For example, a smartphone may employ LTE (Long Term Evolution) for phone calls, WiFi for local area networking, and Bluetooth for headsets. LTE transmissions from the smartphone will cause adjacent channel interference with incoming Bluetooth and WiFi signals. Similarly, Bluetooth and WiFi transmissions from the smartphone will cause adjacent channel interference with incoming LTE signals. This adjacent channel interference can significantly degrade performance not only at the smartphone, but also at connected MWS base stations."}
{"text": "Systems for transporting liquids or gases underground often use a carrier pipe located within a casing (i.e. a second pipe) to protect the carrier pipe from external forces. For example, the pipe may cross a location that has large amounts of heavy vehicle traffic at the surface of the ground, such as where the pipe crosses a road or railroad. Additionally, the ground may undergo significant changes in physical condition such as for example temperature changes or moisture changes that can result in changing external forces being applied to the pipes. Thus, the casing will protect the internal carrier pipe.\nIt is often necessary to install the casing around the carrier pipe after the carrier pipe has been installed. Typically, the retrofit casing must be split in half to form two separate “C-shaped” halves. The halves are then placed around the carrier pipe and then welded back together. Unfortunately, cutting the casing pipe in half is extremely labor intensive. On average, a skilled operator can cut between 6 and 12 inches per minute by hand.\nEmbodiments of the present invention provide improvements over the state of the art."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for manufacturing a molding body having a groove or grooves thereon, on which at least one groove is formed by turning.\n2. Description of the Related Art\nGenerally, there are molding bodies having a predetermined groove or grooves thereon, such as an inserted plate having grooves thereon that is installed inside a mold in order to form a cylindrical lens having grooves formed on a convex surface or a concave surface thereof along its circumferential direction, or to form a prism sheet for use in a backlight device of a liquid crystal display. As a method of manufacturing a molding body having grooves thereon, such as the above-mentioned inserted plate, there has been known a method to rotate a sheet-like workpiece which is wound around a jig in substantially solid cylindrical shape, and a plurality of substantially V-shaped grooves (defined by a shape of an outer edge of a tip) that extend along a circumferential direction of the workpiece are formed substantially in parallel each other by bringing the tip into contact with an outer peripheral face of the workpiece, as disclosed in JP-A-2001-315001.\nIn such a method of manufacturing a molding body having grooves thereon, v-shaped slope faces of the respective grooves are normally formed to be the same in angle. However, recently, it has been requested that these slope faces of the grooves be formed with different angles respectively, as desired. In order to satisfy the above-mentioned request, there have been considered two methods as follows.\nIn one method, a tip C1, having a narrow width and having a rake face C11 with a sharply pointed end, is used as shown in FIGS. 7A and 7B. In this method, by tracing slope faces of virtual grooves Di′ at a front end of the tip C1, the angles of the slope faces of the respective grooves Di can be formed into different angles as desired, as shown in FIG. 7C.\nIn a second method, as shown in FIG. 8A, a tip C2 of which the rake face C21 has a somewhat larger front end angle α (see FIG. 8B), is inclined with the front end thereof fixed as an anchor point so that a base end portion of the rake face C21 approaches a workpiece W. Accordingly, the angle of the slope faces of the groove Di can be changed as desired. In this method, because by inclining the tip C2, the cross-sectional shape of the groove Di at a predetermined position (for example, the front end position of the tip C2) is formed as if the front end shape of the tip C2 at the predetermined position is projected into the workpiece W, as shown in FIGS. 8B and 8C, so that the angle of the slope face of the groove Di can be changed as desired.\nHowever, the first method has a disadvantage that the slope faces of the groove Di are formed by the sharply pointed front end of the tip C1, which brings difficulties in obtaining a good surface roughness. In the second method, a good surface roughness can be obtained since the groove Di is cut by using side edges of the tip C2. However, in order to incline the tip C2 with the front end thereof fixed as the center, the front end of the tip C2 must agree accurately with the rotation center of a jig that rotatably supports the tip C2, and this center-positioning work is extremely tedious. Particularly, in forming a diffractive element, the rotation center of the jig and the front end of the tip C2 have to agree with each other within an error range of 100 nm or less, which is substantially impossible.\nThus, it has been needed to provide a method of manufacturing a molding body having a groove or grooves thereon, which enables slope faces of grooves to be formed with good surface roughness and enables angles of slope faces of grooves to be changed easily."}
{"text": "Yellowish dullness in the skin, meaning that the color tone of the skin becomes to show yellowness with ageing, is simply called “skin yellowness”, hereinafter, unless specified otherwise, and it is considered problematic in the art in terms of skin care. It was once said that the denaturation of proteins in the skin was a major factor of inducing skin yellowness (i.e., saccharification or Maillard reaction); however, nowadays, it is said that the denaturation of proteins in the dermis, existing in a deeper part of the skin (i.e., carbonylation), would be such a major factor.\nSince skin transparency may be deteriorated as the increase of skin yellowness and this may tend to make you look older than your actual age, measures against skin yellowness are distinctly important to keep your youthful appearance. As an actual example of such measures, Patent Literature 1 discloses an agent for inhibiting skin yellowness, which consists of a pumpkin seed extract, while Patent Literatures 2 and 3 disclose external dermal agents which contain a plant extract having an inhibitory effect on skin yellowness. Thus, conventionally known agents, including those which are disclosed in Patent Literatures 1 to 3, are the ones for avoiding skin yellowness by inhibiting denaturation, i.e., saccharification or Maillard reaction, of proteins in the skin.\nPatent Literature 4 discloses a carbonyl-inhibitor for proteins, which contains one or more effective ingredients selected from an olive leaf extract, hydrolyzed pear protein, and lemon extract, and which inhibits neither glycation reaction nor Maillard reaction but inhibits carbonylation so as not to induce yellowness in the skin.\nAs described above, conventional measures against skin yellowness are agents for inhibiting saccharification or Maillard reaction of proteins in the skin so as not to induce yellowness in the skin, meaning that they are so to speak agents for preventing skin yellowness. Thus, it has been made efforts for inhibiting the formation of skin yellowness because it has been recognized in the art that the once-generated skin yellowness could not be diminished to restore the skin to its original condition."}
{"text": "1. Field of the Invention\nThe present invention relates to an image processing apparatus and, more particularly, to pixel conversion processing of an image signal.\n2. Related Background Art\nA conventional image recording/reproduction apparatus comprises a recording unit for recording a video signal picked up by, e.g., a video camera on a recording medium such as a magnetic tape, and a reproduction unit for reproducing the video signal recorded by the recording unit. The video signal reproduced by the reproduction unit is converted in correspondence with a predetermined television system (e.g., NTSC system), and is output to an equipment such as a monitor.\nAs a method of printing a video signal, a video printer is used.\nFor example, when an analog image signal is printed by a video printer, the video printer converts vertically elongated pixels of the NTSC system into square pixels by re-sampling the input analog image signal, stores image data consisting of the converted pixels in its internal frame memory having a predetermined capacity, performs image processing for the image data stored in the frame memory, and forms an image represented by the analog video signal on a transfer medium. The internal frame memory of the video printer has a capacity for, e.g., 640.times.480 pixels.\nWhen a digital video signal output from an image recording/reproduction apparatus which can record/reproduce a digital video signal is to be printed, a video printer having a digital I/F for receiving the digital video signal is used.\nIn this video printer, as in the video printer for receiving the analog video signal, vertically elongated pixels of the input digital video signal are converted into square pixels, the image data size is adjusted to the capacity of its internal frame memory, and the adjusted image data is stored in the frame memory.\nIn this manner, when a conventional image recording/reproduction apparatus outputs a video signal to a video printer, the video printer performs processing such as pixel conversion in correspondence with the characteristics of the video signal. For this reason, the image recording/reproduction apparatus does not execute any pixel conversion processing in correspondence with the characteristics of the video printer as the destination.\nHowever, as described above, since the conventional image recording/reproduction apparatus does not execute any pixel conversion processing in correspondence with the characteristics of a destination equipment such as a video printer, when the conventional image recording/reproduction apparatus is set to be able to execute pixel conversion processing in correspondence with the characteristics of the destination equipment, a means for executing the pixel conversion processing must be added, resulting in an increase in cost."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the fields of molecular biology and virology, and in particular, to the development of gene delivery vehicles. Also disclosed are improved rAAV vector compositions useful in expressing a variety of nucleic acid segments, including those encoding therapeutic proteins polypeptides, peptides, antisense oligonucleotides, and ribozyme constructs, in various gene therapy regimens. Methods are also provided for preparing and using these modified rAAV-based vector constructs in a variety of viral-based gene therapies, and in particular, treatment and prevention of human diseases using conventional gene therapy approaches. The invention also provides rAAV-based vector delivery systems which may be used to assess the relative efficiency and infectivity of a variety of AAV particles having mutations in one or more tyrosine residues of viral capsid proteins.\n2. Description of Related Art\nMajor advances in the field of gene therapy have been achieved by using viruses to deliver therapeutic genetic material. The adeno-associated virus (AAV) has attracted considerable attention as a highly effective viral vector for gene therapy due to its low immunogenicity and ability to effectively transduce non-dividing cells. AAV has been shown to infect a variety of cell and tissue types, and significant progress has been made over the last decade to adapt this viral system for use in human gene therapy.\nIn its normal “wild type” form, recombinant AAV (rAAV) DNA is packaged into the viral capsid as a single stranded molecule about 4600 nucleotides (nt) in length. Following infection of the cell by the virus, the molecular machinery of the cell converts the single DNA strand into a double-stranded form. Only the double-stranded DNA form is useful to the polypeptides of the cell that transcribe the contained gene or genes into RNA.\nAAV has many properties that favor its use as a gene delivery vehicle: 1) the wild type virus is not associated with any pathologic human condition; 2) the recombinant form does not contain native viral coding sequences; and 3) persistent transgenic expression has been observed in many applications.\nThe transduction efficiency of recombinant adeno-associated virus 2 (AAV) vectors varies greatly in different cells and tissues in vitro and in vivo. Systematic studies have been performed to elucidate the fundamental steps in the life cycle of AAV. For example, it has been documented that a cellular protein, FKBP52, phosphorylated at tyrosine residues by epidermal growth factor receptor protein tyrosine kinase (EGFR-PTK), inhibits AAV second-strand DNA synthesis and consequently, transgene expression in virto24, 25 as well as in vivo.19, 27, 28 It has also been demonstrated that EGFR-PTK signaling modulates the ubiquitin/proteasome pathway-mediated intracellular trafficking as well as FKBP52-mediated second-strand DNA synthesis of AAV vectors. In those studies, inhibition of EGFR-PTK signaling led to decreased ubiquitination of AAV capsid proteins, which in turn, facilitated nuclear transport by limiting proteasome-mediated degradation of AAV vectors, implicating EGFR-PTK-mediated phosphorylation of tyrosine residues on AAV capsids."}
{"text": "The present disclosure relates to computer systems and software, and more particularly to techniques for facilitating and automating the provision of services in a cloud environment.\nCloud computing is a model for enabling convenient, on-demand network access to a shared pool of configurable computing resources (e.g., networks, servers, storage, applications, and services). The services provided or accessed through the cloud (or network) are referred to as cloud services. There is a lot of processing that needs to be performed by a cloud service provider to make cloud services available to a subscribing customer. Due to its complexity, much of this processing is still done manually. For example, provisioning resources for providing such cloud services can be a very labor intensive process."}
{"text": "An inductive heating coil that is elastically mounted at a welding gun is disclosed in German Patent Application No. 10 2012 024 708 A1."}
{"text": "The present invention generally relates to the fabrication of water heaters, and more particularly relates to the installation of an externally jacketed foam insulation structure around the storage tank portion of a water heater.\nModern water heaters typically include a storage tank portion adapted to hold a quantity of heated water for on-demand supply to plumbing fixtures such as sinks, tubs, and showers. Projecting outwardly from the storage tank are pipe coupling fittings threadingly connectable to a cold water inlet pipe, a hot water supply pipe, and a temperature/pressure relief pipe. The storage tank is externally insulated by injecting liquid insulating foam into an enclosed insulation space formed between the exterior surface of the storage tank and an external jacket structure disposed outwardly around the storage tank. Piping connections to the storage tank couplings are made through openings formed in the jacket which are spaced outwardly from the couplings in facing relationships with their outer ends.\nDuring the insulating foam injection process it is necessary to prevent the foam from escaping outwardly through the jacket structure pipe connection openings and/or entering the open outer ends of the storage tank pipe coupling fittings. Using conventional methods, this is typically accomplished in two ways. The first method is to provide a fiberglass block having holes formed therethrough to receive the storage tank couplings, position the block against the storage tank exterior with the couplings received in the block holes, and then install the jacket structure. During subsequent injection of the insulating foam the fiberglass block serves as a barrier within the enclosed insulating space to prevent the foam from being forced out through the jacket structure pipe connection openings or entering the storage couplings.\nWhile this is a relatively simple and straightforward approach to forming foam stop barriers, it has two primary disadvantages. First, the fiberglass block must be very carefully sized to sealingly extend between the outer surface of the storage tank and the inner surface of the jacket structure. If even a slight gap exists around the installed block it can easily permit the injected foam to escape from the jacket structure or enter the pipe couplings.\nSecond, the fiberglass block, which tends to be relatively large, typically has a thermal insulation value only about half that of the insulating foam. Accordingly, relative to the foam insulation, the fiberglass block forms a relatively low resistance heat outflow path in the assembled water heater. As energy conservation goals and standards continue to increase, this situation becomes less and less acceptable.\nThe other method conventionally used to form a foam stop barrier around the tank couplings and their associated jacket structure pipe connection openings has been to construct a relatively flat, foam-filled bag having openings therein through which the tank couplings may be extended, taping the bag blanket-like to the tank exterior, and then installing the jacket structure over the bag. Since the bags are filled with foam insulation, they do not present the heat leak problem that the fiberglass blocks do. However, like the fiberglass blocks, the foam filled bags present the potential problem of injected foam leakage past the bags if they are not carefully sized and properly fitted into place within the enclosed insulation space before the foam injection process is initiated. Additionally, the bags are rather tedious and time consuming to fabricate and install, thus undesirably increasing the overall construction cost of the water heater.\nIt can be readily be seen from the foregoing that a need exists for improved foam stop apparatus and methods in the fabrication of externally foam-insulated water heaters. It is accordingly an object of the present invention to provide such improved foam stop apparatus and methods."}
{"text": "1. Field of the Invention\nThis invention relates to the field of computer graphics systems. More particularly, this invention relates to real time generation of complex high resolution stereo 3D images in a head tracked stereo display system.\n2. Art Background\nThe human eye can be modeled as an optical system coupled to a retina, with the retina functioning as a light transducer. The human eye is immersed in physical space filled with light rays. A point source of light exists at every point in the physical space that is a boundary between transparent and opaque surfaces, or a boundary between transparent surfaces of different refractive indices. Human stereo vision is achieved by immersing the two eyes at different locations within the physical space.\nA typical computer graphics system generates images by producing an image plane of light sources, or an image plane of reflectors of light. Such an image plane is approximated by a two dimensional array of point light sources known as pixels. The light produced by the image plane is the sum of the half spheres of light expanding out from each pixel on the image plane. The color and intensity of each pixel is separately controllable, but the distribution of color and intensity across the half sphere usually is not controllable. The pixels that comprise the image plane appear substantially the same when viewed from any angle.\nNevertheless, a head tracked stereo display system simulates the interaction of human eyes with the light rays of the physical space in order to enable a viewer to perceive 3D images. A head tracked stereo display system senses the location in physical space of the viewer's head and eyes, computes a pair of stereo images of a virtual object based upon the location of the head and eyes, and generates the stereo image on a stereo display device.\nA head tracked stereo display system requires that only two images be computed and generated at a time, which results in substantially less computation than holographic systems. Moreover, the pair of stereo images have the same appearance as a hologram. With a head tracked stereo display system, the virtual object appears to remain stationary when the viewer's head tilts, or when the viewer's head moves to look around the side or over the top of the virtual object. For further discussion regarding head tracked stereo display systems, refer to Paley, W. B. Head-tracking Stereo Display, Techniques and Applications, Proceedings of SPIE, Febuary 1992.\nIn a head tracked stereo display system, the scenery and objects of the stereo images must be rendered in a sufficiently short time span to be physiologically acceptable to the human visual system. Past head tracked stereo display systems used commonly available 3D computer graphics image generation hardware and software to produce stereo images with 3D display primitives. However, the performance of such systems is limited by the speed of the 3D graphics rendering hardware. As the complexity of the stereo images increases, the graphics rendering hardware becomes too slow to generate 3D stereo images at an acceptable rate.\nAs will be described, the present method and apparatus generates complex 3D stereo images by dynamically accessing pairs of stero images selected from among a precomputed set of images. The precomputed images are selected based upon a prediction of viewer head location for the next frame of the stereo display.\nThis application is related to application Ser. No. 07/875,042, filed on Apr. 28, 1992, entitled Method and Apparatus for High Resolution Virtual Reality Systems Using Head Tracked Display."}
{"text": "For centuries, since man has been sailing the seas of the world, people have been buried at sea. From depositing a deceased person's body in the water to sprinkling ashes into the ocean, individuals have for centuries desired their final resting place to be in the oceans of the world.\nHowever, when a deceased person's body or ashes are deposited in the sea, there is generally no fixed location that loved ones may visit, as with a conventional cemetery. Because wave currents may disseminate a person's ashes or body into the ocean, loved ones may not have the ability to revisit a point of deposit or other point that may be reserved in memory of the deceased individual.\nIn attempt to accommodate some of these concerns, more recent systems have emerged wherein an individual's ashes may be added to sinkable objects that are dropped into the water at sea. These objects may be constructed such that the ashes are mixed into the materials of the sinkable object such as cement. However, these objects may be dropped at depths and/or locations wherein making a return visit for viewing by boat or even by scuba diving may be difficult, improbable, or even impossible. Plus, the size of such sinkable objects may oftentimes be sufficiently large to overcome any affects of water currents and/or other natural affects, thereby increasing both the maneuvering difficulty during manufacturing and placement at sea size and also the associated costs.\nThus, there is a heretofore unaddressed need to overcome the inefficiencies and shortcomings as described above."}
{"text": "1. Technical Field\nThe present invention relates generally to data compression and decompression and, in particular, to a system and method for compressing and decompressing data based on an actual or expected throughput (bandwidth) of a system that employs data compression. Additionally the present invention relates to the subsequent storage, retrieval, and management of information in data storage devices utilizing either compression and/or accelerated data storage and retrieval bandwidth.\n2. Description of Related Art\nThere are a variety of data compression algorithms that are currently available, both well defined and novel. Many compression algorithms define one or more parameters that can be varied, either dynamically or a-priori, to change the performance characteristics of the algorithm. For example, with a typical dictionary based compression algorithm such as Lempel-Ziv, the size of the dictionary can affect the performance of the algorithm. Indeed, a large dictionary may be employed to yield very good compression ratios but the algorithm may take a long time to execute. If speed were more important than compression ratio, then the algorithm can be limited by selecting a smaller dictionary, thereby obtaining a much faster compression time, but at the possible cost of a lower compression ratio. The desired performance of a compression algorithm and the system in which the data compression is employed, will vary depending on the application.\nThus, one challenge in employing data compression for a given application or system is selecting one or more optimal compression algorithms from the variety of available algorithms. Indeed, the desired balance between speed and efficiency is typically a significant factor that is considered in determining which algorithm to employ for a given set of data. Algorithms that compress particularly well usually take longer to execute whereas algorithms that execute quickly usually do not compress particularly well.\nAccordingly, a system and method that would provide dynamic modification of compression system parameters so as to provide an optimal balance between execution speed of the algorithm (compression rate) and the resulting compression ratio, is highly desirable.\nYet another problem within the current art is data storage and retrieval bandwidth limitations. Modern computers utilize a hierarchy of memory devices. In order to achieve maximum performance levels, modern processors utilize onboard memory and on board cache to obtain high bandwidth access to both program and data. Limitations in process technologies currently prohibit placing a sufficient quantity of onboard memory for most applications. Thus, in order to offer sufficient memory for the operating system(s), application programs, and user data, computers often use various forms of popular off-processor high speed memory including static random access memory (SRAM), synchronous dynamic random access memory (SDRAM), synchronous burst static ram (SBSRAM). Due to the prohibitive cost of the high-speed random access memory, coupled with their power volatility, a third lower level of the hierarchy exists for non-volatile mass storage devices. While mass storage devices offer increased capacity and fairly economical data storage, their data storage and retrieval bandwidth is often much less in relation to the other elements of a computing system.\nComputers systems represent information in a variety of manners. Discrete information such as text and numbers are easily represented in digital data. This type of data representation is known as symbolic digital data. Symbolic digital data is thus an absolute representation of data such as a letter, figure, character, mark, machine code, or drawing.\nContinuous information such as speech, music, audio, images and video, frequently exists in the natural world as analog information. As is well known to those skilled in the art, recent advances in very large scale integration (VLSI) digital computer technology have enabled both discrete and analog information to be represented with digital data. Continuous information represented as digital data is often referred to as diffuse data. Diffuse digital data is thus a representation of data that is of low information density and is typically not easily recognizable to humans in its native form.\nModern computers utilize digital data representation because of its inherent advantages. For example, digital data is more readily processed, stored, and transmitted due to its inherently high noise immunity. In addition, the inclusion of redundancy in digital data representation enables error detection and/or correction. Error detection and/or correction capabilities are dependent upon the amount and type of data redundancy, available error detection and correction processing, and extent of data corruption.\nOne outcome of digital data representation is the continuing need for increased capacity in data processing, storage, and transmittal. This is especially true for diffuse data where increases in fidelity and resolution create exponentially greater quantities of data. Data compression is widely used to reduce the amount of data required to process, transmit, or store a given quantity of information. In general, there are two types of data compression techniques that may be utilized either separately or jointly to encode/decode data: lossless and lossy data compression.\nOver the last decade, computer processor performance has improved by at least a factor of 50. During this same period, magnetic disk storage has only improved by a factor of 5. Thus one additional problem with the existing art is that memory storage devices severely limit the performance of consumer, entertainment, office, workstation, servers, and mainframe computers for all disk and memory intensive operations.\nFor example, magnetic disk mass storage devices currently employed in a variety of home, business, and scientific computing applications suffer from significant seek-time access delays along with profound read/write data rate limitations. Currently the fastest available (15,000) rpm disk drives support only a 40.0 Megabyte per second data rate (MB/sec). This is in stark contrast to the modern Personal Computer's Peripheral Component Interconnect (PCI) Bus's input/output capability of 512 MB/sec and internal local bus capability of 1600 MB/sec.\nAnother problem within the current art is that emergent high performance disk interface standards such as the Small Computer Systems Interface (SCSI-3), iSCSI, Fibre Channel, AT Attachment UltraDMA/100+, Serial Storage Architecture, and Universal Serial Bus offer only higher data transfer rates through intermediate data buffering in random access memory. These interconnect strategies do not address the fundamental problem that all modern magnetic disk storage devices for the personal computer marketplace are still limited by the same typical physical media restriction. In practice, faster disk access data rates are only achieved by the high cost solution of simultaneously accessing multiple disk drives with a technique known within the art as data striping and redundant array of independent disks (RAID).\nRAID systems often afford the user the benefit of increased data bandwidth for data storage and retrieval. By simultaneously accessing two or more disk drives, data bandwidth may be increased at a maximum rate that is linear and directly proportional to the number of disks employed. Thus another problem with modern data storage systems utilizing RAID systems is that a linear increase in data bandwidth requires a proportional number of added disk storage devices.\nAnother problem with most modern mass storage devices is their inherent unreliability. Many modern mass storage devices utilize rotating assemblies and other types of electromechanical components that possess failure rates one or more orders of magnitude higher than equivalent solid state devices. RAID systems employ data redundancy distributed across multiple disks to enhance data storage and retrieval reliability. In the simplest case, data may be explicitly repeated on multiple places on a single disk drive, on multiple places on two or more independent disk drives. More complex techniques are also employed that support various trade-offs between data bandwidth and data reliability.\nStandard types of RAID systems currently available include RAID Levels 0, 1, and 5. The configuration selected depends on the goals to be achieved. Specifically data reliability, data validation, data storage/retrieval bandwidth, and cost all play a role in defining the appropriate RAID data storage solution. RAID level 0 entails pure data striping across multiple disk drives. This increases data bandwidth at best linearly with the number of disk drives utilized. Data reliability and validation capability are decreased. A failure of a single drive results in a complete loss of all data. Thus another problem with RAID systems is that low cost improved bandwidth requires a significant decrease in reliability.\nRAID Level 1 utilizes disk mirroring where data is duplicated on an independent disk subsystem. Validation of data amongst the two independent drives is possible if the data is simultaneously accessed on both disks and subsequently compared. This tends to decrease data bandwidth from even that of a single comparable disk drive. In systems that offer hot swap capability, the failed drive is removed and a replacement drive is inserted. The data on the failed drive is then copied in the background while the entire system continues to operate in a performance degraded but fully operational mode. Once the data rebuild is complete, normal operation resumes. Hence, another problem with RAID systems is the high cost of increased reliability and associated decrease in performance.\nRAID Level 5 employs disk data striping and parity error detection to increase both data bandwidth and reliability simultaneously. A minimum of three disk drives is required for this technique. In the event of a single disk drive failure, that drive may be rebuilt from parity and other data encoded on disk remaining disk drives. In systems that offer hot swap capability, the failed drive is removed and a replacement drive is inserted. The data on the failed drive is then rebuilt in the background while the entire system continues to operate in a performance degraded but fully operational mode. Once the data rebuild is complete, normal operation resumes.\nThus another problem with redundant modern mass storage devices is the degradation of data bandwidth when a storage device fails. Additional problems with bandwidth limitations and reliability similarly occur within the art by all other forms of sequential, pseudo-random, and random access mass storage devices. Typically mass storage devices include magnetic and optical tape, magnetic and optical disks, and various solid-state mass storage devices. It should be noted that the present invention applies to all forms and manners of memory devices including storage devices utilizing magnetic, optical, neural and chemical techniques or any combination thereof.\nYet another problem within the current art is the application and use of various data compression techniques. It is well known within the current art that data compression provides several unique benefits. First, data compression can reduce the time to transmit data by more efficiently utilizing low bandwidth data links. Second, data compression economizes on data storage and allows more information to be stored for a fixed memory size by representing information more efficiently.\nFor purposes of discussion, data compression is canonically divided into lossy and lossless techniques. Lossy data compression techniques provide for an inexact representation of the original uncompressed data such that the decoded (or reconstructed) data differs from the original unencoded/uncompressed data. Lossy data compression is also known as irreversible or noisy compression. Negentropy is defined as the quantity of information in a given set of data. Thus, one obvious advantage of lossy data compression is that the compression ratios can be larger than that dictated by the negentropy limit, all at the expense of information content. Many lossy data compression techniques seek to exploit various traits within the human senses to eliminate otherwise imperceptible data. For example, lossy data compression of visual imagery might seek to delete information content in excess of the display resolution or contrast ratio of the target display device.\nOn the other hand, lossless data compression techniques provide an exact representation of the original uncompressed data. Simply stated, the decoded (or reconstructed) data is identical to the original unencoded/uncompressed data. Lossless data compression is also known as reversible or noiseless compression. Thus, lossless data compression has, as its current limit, a minimum representation defined by the entropy of a given data set.\nA rich and highly diverse set of lossless data compression and decompression algorithms exist within the current art. These range from the simplest “adhoc” approaches to highly sophisticated formalized techniques that span the sciences of information theory, statistics, and artificial intelligence. One fundamental problem with almost all modern approaches is the compression ratio to encoding and decoding speed achieved. As previously stated, the current theoretical limit for data compression is the entropy limit of the data set to be encoded. However, in practice, many factors actually limit the compression ratio achieved. Most modern compression algorithms are highly content dependent. Content dependency exceeds the actual statistics of individual elements and often includes a variety of other factors including their spatial location within the data set.\nOf popular compression techniques, arithmetic coding possesses the highest degree of algorithmic effectiveness, and as expected, is the slowest to execute. This is followed in turn by dictionary compression, Huffman coding, and run-length coding with respectively decreasing execute times. What is not apparent from these algorithms, that is also one major deficiency within the current art, is knowledge of their algorithmic efficiency. More specifically, given a compression ratio that is within the effectiveness of multiple algorithms, the question arises as their corresponding efficiency.\nWithin the current art there also presently exists a strong inverse relationship between achieving the maximum (current) theoretical compression ratio, which we define as algorithmic effectiveness, and requisite processing time. For a given single algorithm the effectiveness over a broad class of data sets including text, graphics, databases, and executable object code is highly dependent upon the processing effort applied. Given a baseline data set, processor operating speed and target architecture, along with its associated supporting memory and peripheral set, we define algorithmic efficiency as the time required to achieve a given compression ratio. Algorithmic efficiency assumes that a given algorithm is implemented in an optimum object code representation executing from the optimum places in memory. This is almost never achieved in practice due to limitations within modern optimizing software compilers. It should be further noted that an optimum algorithmic implementation for a given input data set may not be optimum for a different data set. Much work remains in developing a comprehensive set of metrics for measuring data compression algorithmic performance, however for present purposes the previously defined terms of algorithmic effectiveness and efficiency should suffice.\nVarious solutions to this problem of optimizing algorithmic implementation are found in U.S. Pat. Nos. 6,195,024 and 6,309,424, issued on Feb. 27, 2001 and Oct. 30, 2001, respectively, to James Fallon, both of which are entitled “Content Independent Data Compression Method and System,” and are incorporated herein by reference. These patents describe data compression methods that provide content-independent data compression, wherein an optimal compression ratio for an encoded stream can be achieved regardless of the data content of the input data stream. As more fully described in the above incorporated patents, a data compression protocol comprises applying an input data stream to each of a plurality of different encoders to, in effect, generate a plurality of encoded data streams. The plurality of encoders are preferably selected based on their ability to effectively encode different types of input data. The final compressed data stream is generated by selectively combining blocks of the compressed streams output from the plurality of encoders based on one or more factors such as the optimal compression ratios obtained by the plurality of decoders. The resulting compressed output stream can achieve the greatest possible compression, preferably in real-time, regardless of the data content.\nYet another problem within the current art relates to data management and the use of existing file management systems. Present computer operating systems utilize file management systems to store and retrieve information in a uniform, easily identifiable, format. Files are collections of executable programs and/or various data objects. Files occur in a wide variety of lengths and must be stored within a data storage device. Most storage devices, and in particular, mass storage devices, work most efficiently with specific quantities of data. For example, modern magnetic disks are often divided into cylinders, heads and sectors. This breakout arises from legacy electromechanical considerations with the format of an individual sector often some binary multiple of bytes (512, 1024, . . . ). A fixed or variable quantity of sectors housed on an individual track. The number of sectors permitted on a single track is limited by the number of reliable flux reversals that can be encoded on the storage media per linear inch, often referred to as linear bit density. In disk drives with multiple heads and disk media, a single cylinder is comprised of multiple tracks.\nA file allocation table is often used to organize both used and unused space on a mass storage device. Since a file often comprises more than one sector of data, and individual sectors or contiguous strings of sectors may be widely dispersed over multiple tracks and cylinders, a file allocation table provides a methodology of retrieving a file or portion thereof. File allocation tables are usually comprised of strings of pointers or indices that identify where various portions of a file are stored.\nIn-order to provide greater flexibility in the management of disk storage at the media side of the interface, logical block addresses have been substituted for legacy cylinder, head, sector addressing. This permits the individual disk to optimize its mapping from the logical address space to the physical sectors on the disk drive. Advantages with this technique include faster disk accesses by allowing the disk manufacturer greater flexibility in managing data interleaves and other high-speed access techniques. In addition, the replacement of bad media sectors can take place at the physical level and need not be the concern of the file allocation table or host computer. Furthermore, these bad sector replacement maps are definable on a disk by disk basis.\nPractical limitations in the size of the data required to both represent and process an individual data block address, along with the size of individual data blocks, governs the type of file allocation tables currently in use. For example, a 4096 byte logical block size (8 sectors) employed with 32 bit logical block addresses. This yields an addressable data space of 17.59 Terabytes. Smaller logical blocks permit more efficient use of disk space. Larger logical blocks support a larger addressable data space. Thus one limitation within the current art is that disk file allocation tables and associated file management systems are a compromise between efficient data storage, access speed, and addressable data space.\nData in a computer has various levels of information content. Even within a single file, many data types and formats are utilized. Each data representation has specific meaning and each may hold differing quantities of information. Within the current art, computers process data in a native, uncompressed, format. Thus compressed data must often be decompressed prior to performing various data processing functions or operations. Modern file systems have been designed to work with data in its native format. Thus another significant problem within the current art is that file systems are not able to randomly access compressed data in an efficient manner.\nFurther aggravating this problem is the fact that when data is decompressed, processed and recompressed it may not fit back into its original disk space, causing disk fragmentation or complex disk space reallocation requirements. Several solutions exist within the current art including file by file and block structured compressed data management.\nIn file by file compression, each file is compressed when stored on disk and decompressed when retrieved. For very small files this technique is often adequate, however for larger files the compression and decompression times are too slow, resulting in inadequate system level performance. In addition, the ability to access randomly access data within a specific file is lost. The one advantage to file by file compression techniques is that they are easy to develop and are compatible with existing file systems. Thus file by file compressed data management is not an adequate solution.\nBlock structured disk compression operates by compressing and decompressing fixed block sizes of data. Block sizes are often fixed, but may be variable in size. A single file usually is comprised of multiple blocks, however a file may be so small as to fit within a single block. Blocks are grouped together and stored in one or more disk sectors as a group of Blocks (GOBs). A group of blocks is compressed and decompressed as a unit, thus there exists practical limitations on the size of GOBs. Most compression algorithms achieve a higher level of algorithmic effectiveness when operating on larger quantities of data. Restated, the larger the quantity of data processed with a uniform information density, the higher the compressions ratio achieved. If GOBs are small compression ratios are low and processing time short. Conversely, when GOBS are large compression ratios are higher and processing time is longer. Large GOBs tend to perform in a manner analogous to file by file compression. The two obvious benefits to block structured disk compression are psuedo-random data access and reduced data compression/decompression processing time.\nSeveral problems exist within the current art for the management of compressed blocks. One method for storage of compressed files on disk is by contiguously storing all GOBs corresponding to a single file. However as files are processed within the computers, files may grow or shrink in size. Inefficient disk storage results when a substantial file size reduction occurs. Conversely when a file grows substantially, the additional space required to store the data may not be available contiguously. The result of this process is substantial disk fragmentation and slower access times.\nAn alternate method is to map compressed GOBs into the next logical free space on the disk. One problem with this method is that average file access times are substantially increased by this technique due to the random data storage. Peak access delays may be reduced since the statistics behave with a more uniform white spectral density, however this is not guaranteed.\nA further layer of complexity is encountered when compressed information is to be managed on more than one data storage device. Competing requirements of data access bandwidth, data reliability/redundancy, and efficiency of storage space are encountered.\nThese and other limitations within the current art are solved with the present invention."}
{"text": "The present invention relates generally to spacecraft engineering and design. More particularly, the present invention relates to structural, fluid, and electrical connectors for interconnection of spacecraft modules. Furthermore, the present invention relates to spacecraft docking and separation systems for providing structural, fluid, and electrical connection between components.\nThere is an enormous benefit to docking or connecting spacecraft and spacecraft components. The many applications include launching a spacecraft for rendezvous in orbit with another spacecraft. Furthermore, the ability to connect or disconnect spacecraft modules has application where a broken module can be jettisoned and replaced. Electronic and structural components are also connected prior to launch to form a spacecraft.\nStructural, fluid and electrical connectors are often combined to create an “interface” connection which is employed in a wide variety of applications. Current spacecraft and launch vehicle interfaces employ various electrical connectors and separation systems including pyrotechnically actuated clamp bands, separation nuts, and separation bolts. The interface must be capable of transferring loads between the two structures. The loads may include vibration, acceleration, thermal conduction, and static loads. For example, U.S. Patent Application Publication No. 2002/0164204 describes a spacecraft interface structure for mechanically connecting various spacecraft components including a spacecraft to a launch vehicle.\nStructural, fluid and electrical interface connectors have been developed for in-space docking and separation. For example, U.S. Pat. No. 8,006,937 describes a docking interface in which alignment cups are heated to an austenitic temperature to form a mechanical connection. Power and fluids may be transferred between the two spacecrafts. Fluid, such as a propulsion fluid, is transferred through a central conduit, and data and power can be transferred through the mechanical coupling. Similarly, U.S. Patent Application Publication No. 2012/0000575 describes a docking interface for refueling satellites in space. The docking assembly includes a connector including male and female components for providing fluid transfer. The connector provides electrical power and data transfer as well.\nPast spacecraft have been custom designed and optimized for a particular mission or payload wherein electronic components are mounted inside box-like modules. The spacecraft electronic components provide a single spacecraft function, such as power, thermal, structural, telemetry, tracking, control or processing, or a specific payload function, such as communications or surveillance. The modules are mounted to a spacecraft frame and interconnected with extremely complex wiring harnesses. Unfortunately, the designs cannot be adapted to accomplish different missions as changing any part or sub-system typically has required significant rework of the design.\nModular spacecraft are attempts to transition away from custom designed spacecraft toward multi-use designs and mass production in an effort to reduce the cost of spacecraft development. The modular spacecraft include a bus which provides a general purpose spacecraft platform. Various payloads can be mounted to the spacecraft platform utilizing a standard structural, fuel and electric interface. Again, the interface connector is of paramount importance.\nEven more recently, spacecraft have been designed by combing homogeneous cells. Each cell possesses the traditional architecture of a spacecraft including structure, power, fuel, attitude control and determination, satellite processing, etc. Each cell is substantially identical so as to be manufactured inexpensively and quickly. These cells are combined to create larger and larger platforms to support payload functions such as communications and surveillance. Though there is substantial redundancy by each cell incorporating all spacecraft sub-system capabilities, the extra costs are more than made up for by mass production savings and rapid assembly. Of course, the satellite reconfigurable cell concept also requires that cells be interconnected by a structural, fluid and electrical interface.\nUnfortunately, previous structural, fluid and electrical interface connectors have suffered from various drawbacks.\nThus, there is a need for an improved structural, fluid and electrical connector.\nThere is also a need for an improved spacecraft connector which can be docked and separated in space.\nAdditionally, it would be advantageous to provide a connector that can be utilized as a launch lock for affixing a spacecraft to a launch vehicle.\nAlso desirable would be a spacecraft connector that can be utilized to connect and function as an interface between homogeneous cells of a spacecraft.\nFurthermore, it would be desirable to provide a spacecraft connector which is capable of functioning as a propulsion thruster when not functioning as a spacecraft interface connector."}
{"text": "U.S. Pat. No. 5,096,950 discloses the co-use of certain NOR (N-alkoxy) hindered amines with a brominated Sb2O3-containing flame retardant in polypropylene.\nU.S. Pat. No. 5,393,812 discloses polyolefin compositions which are made flame retardant by a combination of a halogenated hydrocarbyl phosphate or phosphonate ester flame retardant in combination with a alkoxyamine functional hindered amine.\nU.S. Pat. No. 5,844,026 discloses polyolefin compositions comprising certain NOR hindered amines and certain traditional flame retardants.\nU.S. Pat. No. 6,117,995 discloses that certain N-alkoxy hindered amines may be used as flame retardants for organic polymers.\nU.S. Pat. No. 6,271,377 discloses polyolefin compositions that comprise N-hydroxyalkoxy hindered amines and a halogenated flame retardant.\nU.S. Pat. No. 6,309,987 and equivalent WO 99/54530 teach polyolefin non-woven flame retardant fabrics that comprise N-alkoxyamines.\nA Revolutionary UV Stable Flame Retardant System for Polyolefins—R. Srinivasan, A. Gupta and D. Horsey, Int. Conf. Addit. Polyolefins 1998, 69-83, teaches polyolefins comprising certain NOR hindered amines with halogen and phosphorus containing traditional flame retardants.\nAdvances in a Revolutionary Flame Retardant System for Polyolefins—R. Srinivasan, B. Rotzinger, Polyolefins 2000, Int Conf Polyolefins 2000, 571-581, teaches polyolefins comprising certain NOR hindered amines with brominated and phosphorus containing flame retardants.\nEP 0792911 A2, discloses polyolefin compositions that comprise alkoxyamine functional hindered amines and tris(trihalogenopentyl) phosphate flame retardants.\nWO 99/00450, copending U.S. application Ser. Nos. 09/502,239, filed Nov. 3, 1999, and 09/714,717, filed Nov.16, 2000, disclose the use of certain N-alkoxy hindered amines as flame retardants.\nEP 1104766 discloses cross-linked phenoxyphosphazene compounds as flame retardants for synthetic resins.\nThe flame retardant (FR) market today is comprised of products which function to interfere with the combustion process by chemical and/or physical means. Mechanistically these FRs have been proposed to function during combustion of an article in either the gas phase, the condensed phase or both. The organohalogens are proposed to generate halogen species (e.g. HX) which interferes in the gas phase with free radical organic “fuel” from the polymer substrate. Synergists are proposed to react with HX to form additional chemical species with interfere with combustion in the gas phase, such as reaction of antimony oxide with HX to form antimony halide and water vapor. Antimony compounds such as antimony trioxide also act as a radical scavenger forming antimony halides. Thus, it can inhibit the propagation of the fire.\nAlthough antimony compounds are efficient in terms of cost performance, it recently raised a lot of concern because of the toxicity of the byproducts which are formed during combustion in the presence of a halogenated flame retardant. Antimony oxides often contain trace amounts of arsenic compounds which are suspected carcinogens. Because of these ecological concerns, there is a motion to replace antimony trioxide in the present commercial flame retardant applications. However, it is very difficult to find an effective synergist which is both enviromentally friendly and efficient as far as the cost performance is concerned.\nAnother reason to add flame retardant additives is to prevent dripping during the application of the fire. Dripping during combustion is the process of the separation of parts of the polymer from the matrix in the shape of droplets. Most often, the droplets are flaming and are imposing tremendous danger for fire spread. It is a common measure to add fillers such talc in large amounts to the polymer, with some negative consequences on the mechanical properties. Other fillers sometimes used include calcium carbonate, magnesium carbonate, zinc borate, silicates, silicones, glass fibres, glass bulbs, asbestos, kaolin, mica, barium sulfate, calcium sulfate, metal oxides, hydrates and hydroxides such as zinc oxide, magnesium hydroxide, alumina trihydrate, silica, calcium silicate, magnesium silicate.\nIt has been found that polymers with good flame retardant properties are prepared when a sterically hindered amine of the nitroxyl, hydroxylamine or alkoxyamine class is added thereto together with a phosphazene flame retardant. With this combination, antimony compounds and fillers may be largely reduced or replaced. As the instant nitroxyl, hydroxylamine and alkoxyamine compounds are active as stabilizers, the polymer compositions of the invention are efficiently protected from the deleterious effects of light, oxygen and/or heat."}
{"text": "HVACR systems that operate using a vapor-compression cycle generally comprise a compressor, a condenser, an expansion valve and an evaporator interconnected by a line having an interior within which a refrigerant fluid is circulated. Although this technology is mature and is currently used in a wide variety of commercial and residential applications, malfunctions can arise as a result of leaks in the system or failure of one or more components. In many instances, the malfunction may not be catastrophic, e.g. wherein the system ceases to heat or cool altogether, but may result in a gradual or progressive decrease in performance and/or efficiency that can nevertheless create spoilage of foodstuffs and other inventory, for example, or other problems.\nPeriodic maintenance of HVACR systems can be time consuming, expensive and in many cases unnecessary at the time performed. Recognizing that a failure to maintain such systems will eventually cause a problem, the issue becomes how often such maintenance should be performed and what should be done. Approaches that rely solely on the passage of time are often ineffective and ignore the specific requirements of a particular installation and/or type of system.\nAutomated preventative maintenance devices for HVACR systems have been proposed, such as disclosed, for example, in U.S. Pat. No. 5,596,507. The device disclosed in the '507 patent relies upon a number of temperature sensors and electrical current sensors to detect a variety of operating parameters of the system, and to provide inputs to a computer capable of analyzing the data and identifying potential trouble spots in the system that may need maintenance. While systems of this type may be effective, they are not economically feasible for residential applications and many smaller commercial operations."}
{"text": "A variety of containers such as buckets, whether prepackaged with a material or readily empty for use, often include a carrying handle that comprises a simple wire bail extending from one side of the bucket to the other. Opposite ends of the wire bail are typically bent at a right angle to provide a short segment that is securely retained within an aperture provided for on opposite sides of the bucket allowing the handle to pivot about a horizontal axis. The wire bail handle provides limited comfort when carrying a heavy bucket, especially for extended periods of time. The wire bail handle can be uncomfortable to hold, particularly when the bucket is loaded with heavy material such as mortar, rocks and dirt. Containers incorporating the use of wire bail handles can be found in local hardware or retail stores. For example, there are buckets sized and designed to hold paint, sealants, acrylic coating, tar pitch, dry cement, and other liquids or solids. The buckets come in a variety of different shapes, and are often sized to hold one to five gallons of a liquid material or sized to hold a certain amount of weight. Wire bail handles are inexpensive to fabricate, easy to assemble on buckets, and provide a quick and easy method for carrying or transporting a bucket.\nIn an effort to provide comfort to a person's hands and fingers when carrying a heavy bucket having a wire bail handle, industry has developed a small, plastic cylindrical tube that is generally located in a central holding region of the wire bail handle. The plastic tube typically includes a longitudinal opening extending along the center of the tube for the wire bail to extend through the opening. Because the wire bail is curved on opposite sides of the plastic tube, the cylindrical tube is prevented from sliding back and forth on the wire bail. Though the use of a plastic tube provides some improvement, it does not resolve the issue of having discomfort to the hands and fingers when carrying a heavy bucket. The plastic cylindrical tube is too rigid and creates pain, difficulty, and often leaves physical marks on a person's hands or fingers when carrying a heavy bucket. For example, when carrying a loaded bucket, the rigid tube compresses against the underside of the fingers displacing the skin and rubbing firmly on the bones of the fingers.\nThere have been several attempts to provide solutions to the problem stated above in the prior art. Other than completely replacing the prior art bucket handle, some of the most popular designs provide conventional prior art devices that include a cushion or pad-like material used for covering the handle. For example, in one conventional design, a number of inter-fitting pieces are installed on an already existing handle of a bucket to help enhance the conformability of the handle. The inter-fitting pieces may comprise hard rubber, soft plastic or foam-like elements that are secured together by fasteners such as screws or adhered together using glue or a strong, bonding adhesive. Another exemplary device provides an air cushion grip that is also fixed around an already existing handle. The cushion grip comprises an outer flat layer and a bottom layer provided with a plurality of air cells. The air cells are inflated or deflated, with the use of an air valve or pump, for adjusting the elasticity and shock-absorbing ability of the grip. Still another example of a conventional prior art device includes a pad or wrap configured for wrapping around or being attached to the outer surface of a handle. The pad generally comprises a square, laminate material having an open or closed-cell foam cushion and fasteners disposed along the outer surface of the laminate for securely wrapping pad around the handle. The fasteners often include snaps or straps. One drawback of such devices is that the layout of the fasteners results in either a pad rotating loosely around the handle because the pad cannot be attached firmly onto the handle, or that the underlying fasteners result in a bulky, lesser comfortable handle. Some prior art devices have also developed the use of non-slip surfaces or have include friction methods to prevent the pad from rotating around the handle. Still, other fasteners have been incorporated to ease the manner of wrapping a pad around a handle and to provide for a firm attachment. Such devices have incorporated the use of hooks and loops or Velcro™ fasteners that are disposed along the edges or central region of the pad to allow an easy and firm method of wrapping a cushion or pad around a handle. Though this conventional handle wrap or padding provides some level of comfort, it presents certain drawbacks. Over time, the hooks grab onto materials such as lint, dust, or other debris thereby reducing the bonding strength of the hook and loop from repeated use. Further, the loops can become extended and even break as a result of continued use.\nAs such, the above cited solutions provide little resolve to alleviating discomfort and pain to a user's hands when carrying a heavy bucket by the handle. What is desired is a handle wrap or padding that is easily retrofitted onto an existing handle, attaches securely to a handle without the need for fasteners, and more importantly includes shock absorbing characteristics, and relieves pressure to the hands and fingers of a user when carrying a bucket or other container loaded with heavy material such as mortar, paint, powder, dirt, rocks, etc. . . ."}
{"text": "A socket is used when inspecting a PCBA (Printed Circuit Board Assembly) mounted with electronic components, and this socket has inspection probes that make electrical connections to the PCBA. For example, it is possible to use a socket that is capable of contacting the probes to connection terminals provided on a lower surface of the PCBA, from under the lower surface of the PCBA, as proposed in Japanese Laid-Open Patent Publication No. 2011-123015, for example.\nHowever, in a case in which the PCBA has a structure such that the probes cannot contact the connection terminals provided on the lower surface of the PCBA, from under the lower surface of the PCBA, a special structure for enabling inspection needs to be provided on the PCBA. For example, the special structure for enabling inspection may be an inspection pattern that is provided on the lower surface of the PCBA and contactable by the probes upon inspection, or an inspection pattern that is provided in an extended region of the PCBA extended in a horizontal direction of the PCBA and contactable by the probes upon inspection. But the special structure for enabling inspection increases a size of the PCBA, increases a cost of the PCBA, or increases both the size and the cost of the PCBA."}
{"text": "Over the years, general purpose railway boxcars have progressed from relatively simple wooden structures mounted on flat cars to more elaborate arrangements including insulated walls and refrigeration equipment. Various types of insulated railway boxcars are presently manufactured and used. A typical insulated railway boxcar includes an enclosed structure mounted on a railway car underframe. The enclosed structure generally has an outer shell, one or more layers of insulation and interior paneling. The outer shell of such railway boxcars often has an exterior surface formed from various types of metal such as steel or aluminum. The interior paneling is often formed from wood and/or metal as desired for the specific application. For some applications the interior paneling has been formed from fiber reinforced plastic (FRP). Various types of sliding doors including plug type doors are generally provided on each side of conventional railway boxcars for loading and unloading freight. Conventional railway boxcars are assembled from various pieces of wood, steel and/or sheets of composite materials such as fiberglass reinforced plastic. Significant amounts of raw material, labor and time are often required to complete the manufacture and assembly of conventional railway boxcars.\nThe underframe for many railway boxcars include a center sill with a pair of end sills and a pair of side sills arranged in a generally rectangular configuration corresponding approximately with the dimensions for the floor of the railway boxcar. Cross bearers and cross ties are provided to establish the desired rigidity and strength for transmission of vertical loads from the side sills to the center sill and for dissipating horizontal end loads on the center sill to other portions of the underframe. A plurality of longitudinal stringers are also provided on each side of the center sill to support the floor of the enclosed structure. Examples of such railway car underframes are shown in U.S. Pat. Nos. 2,783,718 and 3,266,441.\nFor many years various techniques have been used to build fiberglass boat hulls. Many of these hulls have been fabricated using wet layup techniques in which each layer of material such as fiberglass or carbon fiber is first wetted with the desired resin such as polyester or vinylester and then laid in an open mold. Recently, vacuum bagging techniques have been combined with wet layup techniques to control the emission of volatile organic compounds. Vacuum bagging also produces a stronger structure by eliminating air pockets and excess resin in the finished product.\nMore recently, vacuum bagging techniques have been combined with an enhanced resin delivery system which allows the use of a closed molding system and dry layup of core layers and fiber reinforcing layers such as fiberglass in the mold. This process may sometimes be referred to as composite resin infusion molding. U.S. Pat. Nos. 4,902,215; 5,052,906 and 5,316,462 provide additional information concerning this type of vacuum bagging process to form a fiberglass reinforced composite article.\nVarious types of load dividers and freight securing systems have previously been used to prevent undesired movement of freight contained within a railway boxcar. The use of such systems is particularly important when a railway boxcar is only partially loaded. Examples of such systems are shown in U.S. Pat. No. 5,370,482 entitled \"Cargo Securement System\" and U.S. Pat. No. 5,386,674 entitled \"Two Piece Bulkhead Door for Rail Cars and the Like.\" All patents noted in the Background of the Invention are incorporated by reference for all purposes within this application."}
{"text": "1. Field of the Invention\nThis invention relates to a sheet supply system, and more particularly to a sheet supply system which feeds out sheets stacked on a sheet support table one by one from the sheet support table.\n2. Description of the Related Art\nFor example, in a stencil printer, a number of printing papers are supplied to a printing system one by one without lost feed (referring to malfunction that no printing paper is supplied) or multiple feed (referring to malfunction that a plurality of printing paper are supplied at one time).\nIn the stencil printer, a number of copies are often printed from one original, and accordingly a number of printing papers are stacked on a paper support table. In such a case, it is preferred that the uppermost printing paper in the stack on the support table be constantly positioned at a predetermined level (will be referred to as \"the paper feed-out level\", hereinbelow) so that a paper feed-out system can act on the uppermost printing paper constantly in an optimal manner. For this purpose, the paper support table is moved upward as the level of the uppermost printing paper in the stack remaining on the paper support table lowers.\nFIG. 8 shows an example of a conventional paper supply system for a stencil printer. In FIG. 8, reference numeral 51 denotes a paper support table on which a stack of printing papers P is supported. The paper support table 51 is moved up and down by a table drive means (not shown). The printing papers P are fed out from the paper support table 51 one by one by a paper feed-out unit 53. The paper feed-out unit 53 comprises a pair of paper feed rollers (a scraper roller 54a and a pick-up roller 54b) which are mounted for rotation on one bracket and are rotated in synchronization with each other. The paper feed-out unit 53 is rotatable about a shaft 53a and is moved upward when the paper support table 51 is moved upward and the uppermost printing paper in the stack of the printing papers P pushes upward the scraper roller 54a. The shaft 53a is urged by a spring means (not shown) so that the contact pressure between the scraper roller 54a and the uppermost printing paper increases as the paper feed-out unit 53 is lifted higher.\nA limit switch 52, which may be, for instance, a microswitch or a photo-interrupter, is disposed to be opened when the paper feed-out unit 53 is moved upward to a predetermined level. That is, when paper support table 51 is moved upward to push the paper feed-out unit 53 upward and the limit switch 52 is opened, the table drive means is de-energized and the paper support table 51 is stopped at the level. When the printing papers P on the support table 51 are fed out from the paper support table 51 and the level of the uppermost printing paper lowers, the limit switch 52 is closed again and the table drive means is energized to move upward the paper support table 51. In this manner, the uppermost printing paper is constantly kept at the level at which the limit switch 52 is opened, thereby keeping constant the contact pressure between the uppermost printing paper and the scraper roller 54a during paper supply operation of the paper supply system. The level of the uppermost printing paper at which the limit switch 52 is opened is the aforesaid \"paper feed-out level\".\nHowever the conventional system is disadvantageous in that since the paper feed-out level is controlled by opening and closure of the limit switch 52, the accuracy in positioning the uppermost printing paper at the paper feed-out level is governed by the sensitivity of the limit switch 52, which fluctuates from switch to switch and changes with time. When the paper feed-out level fluctuates, the contact pressure between the uppermost printing paper and the scraper roller 54a fluctuates, which causes the aforesaid lost feed or the aforesaid multiple feed. Further, it takes a movement of the paper support table 51 through a certain distance to open and close the limit switch 52. The distance is generally larger than the thickness of one printing paper and is generally equal to triple of the thickness of the printing paper. Accordingly, the level of the uppermost printing paper is adjusted only once per three printing papers, which results in fluctuation in the contact pressure during supply of the three printing papers.\nThere are used various types of printing papers which are different in thickness, quality and the like. The contact pressure optimal to feed the printing papers one by one differs according to the thickness, quality and the like of the printing paper. Thus there has been proposed, as disclosed in Japanese Unexamined Patent Publication No. 8(1996)-259099, a paper supply system in which the contact pressure between the uppermost printing paper and the scraper roller (i.e., the paper feed-out level) can be changed. In the paper supply system, the position of the limit switch 52 relative to the paper feed-out unit 53 is changed in two steps by operation of a lever so that the contact pressure between the uppermost printing paper and the scraper roller 54a at the level of the uppermost printing paper at which the limit switch 52 is opened is changed in two steps. When thick printing papers are to be supplied, the contact pressure is increased and when thin printing papers are to be supplied, the contact pressure is reduced.\nAlso this system cannot be free from the aforesaid drawbacks inherent to the limit switch and since the paper feed-out level is changed by changing the position of the limit switch relative to the paper feed-out unit, an additional mechanism is required, which adds to the manufacturing cost of the system.\nFurther printing papers of different thicknesses can be sometimes stacked on the paper support table, and accordingly there is a demand for a paper supply system which can automatically change the paper feed-out level according to the thickness of the printing paper."}
{"text": "1. Field of the Invention\nThe present invention relates to a safety switch that ensures the bi-metallic plate to be deformed as desired when overloaded.\n2. The Prior Arts\nA conventional switch device, especially for those switches using bi-metallic plate to prevent from being burn when an overload happens, generally includes a bi-metallic plate which is deformed when overloaded so as to separate the two contact points respectively located on the bimetallic plate and one of the two terminals. Some inherent shortcomings for these conventional safety switches are experienced. There are too many parts involved in the safety switches and a longer period of time is required when assembling the switches, so this increases the cost of the products. The parts might be arranged inaccurately and affect the deformation of the bi-metallic plate. Once the bimetallic plate is deformed to cut off the circuit, because of the improper arrangement of the parts as mentioned above, the bi-metallic plate could deform to re-connect the two contact points to connect the circuit again. Because of the inaccuracy of the deformation of the bi-metallic plate, the switch member cannot be set at “OFF” position after the bimetallic plate is deformed to cut off the circuit.\nTherefore, it is desired to have a safety switch that allows the bi-metallic plate to deform toward a desired direction and has enough space for the deformation of the bi-metallic plate to prevent the bimetallic plate from bouncing back to connect the circuit again."}
{"text": "1. Field of the Invention\nThe present invention relates to a variable displacement compressor, and more specifically, to a variable displacement compressor with an improved structure for suppressing axial vibration of a drive shaft.\n2. Description of Related Art\nVariable displacement compressors, particularly, variable displacement inclined plate compressors, are known in the art. For example, a known structure of a variable displacement inclined plate compressor is constructed as depicted in FIG. 2. In FIG. 2, variable displacement inclined plate compressor 1xe2x80x2 has cylindrical housing 2, front end plate 3 closing the front end of housing 2, and cylinder head 4 closing the rear end of housing 2. Valve plate 5 is interposed between the rear end of housing 2 and cylinder head 4. Cylinder block 6 is disposed in the rear half of housing 2. A plurality of cylinder bores 7, center bore 8, and communication path 9, are formed in cylinder block 6. Cylinder bores 7 are arranged radially about a center axis of the compressor at even radial intervals. Cylinder bores 7 extend in a front/rear direction. The front half of housing 2 forms crank chamber 10.\nDrive shaft 11 extends into housing 2 in the front/rear direction through front end plate 3. The rear end portion of drive shaft 11 is inserted into center bore 8. Drive shaft 11 is rotatably supported by front end plate 3 via radial bearing 12 and by cylinder block 6 via radial bearing 13. The rear end of drive shaft 11 is axially supported by thrust bearing 14 and adjusting screw 15xe2x80x2 threaded into center bore 8. Axial gaps of thrust bearing 14 and thrust bearing 17, which are described in detail later, are adjusted by controlling the degree to which adjusting screw 15xe2x80x2 is threaded into center bore 8, which thereby, adjusts the axial loads applied to thrust bearings 14 and 17. Lip seal 34 is disposed at a front side of radial bearing 12.\nRotor 16 is disposed in crank chamber 10 and fixed to drive shaft 11. Rotor 16 is supported in the axial direction of drive shaft 11 by front end plate 3 via thrust bearing 17. Arm portion 16a is formed by the rear end portion of rotor 16. Slot 16b is defined in arm portion 16a. Spherical bush 18 is slidably fitted onto drive shaft 11 at a rear position of rotor 16, in the axial direction of drive shaft 11. Coil spring 19 is interposed between rotor 16 and spherical bush 18. Disc-type inclined plate 20 is provided slidably and rotatably on spherical bush 18.\nArm portion 20a is provided on one side of inclined plate 20. Arm portion 20a extends toward selvage portion 16a of rotor 16. Arm portion 20a has hole 20b defined at a position corresponding to slot 16b. Pin 21 is inserted into slot 16b and hole 20b to connect rotor 16 and inclined plate 20 while allowing the variable inclination of inclined plate 20.\nA pair of semi-spherical sliding shoes 22 are provided slidably on both surfaces of a radial outer portion of inclined plate 20. A plurality of pairs of semi-spherical sliding shoes 22 are disposed radially about inclined plate 20 at even intervals. Each pair of semi-spherical sliding shoes 22 are held slidably in each piston rod 23. Each piston rod 23 extends into a corresponding cylinder bore 7 in the rear direction, and forms a piston 24 slidably inserted into the corresponding cylinder bore 7.\nSuction port 25 and discharge port 26 are provided on valve plate 5 in correspondence with each cylinder bore 7. A suction valve (not shown) and a discharge valve (not shown) are provided for controlling the flow of fluid through suction port 25 and discharge port 26, respectively. Suction chamber 27, communicating with suction port 25, and discharge chamber 28, communicating with discharge port 26, are formed in cylinder head 4. Suction chamber 27 communicates with inlet port 29. Discharge chamber 28 communicates with outlet port 30.\nCommunication path 9, formed in cylinder block 6, communicates with crank chamber 10 and suction chamber 27. Bellows 31 is disposed in communication path 9. Electromagnetic clutch 32 is provided at a front position of front end plate 3.\nIn variable displacement inclined plate compressor 1xe2x80x2, an external driving force is transmitted from an external drive source (not shown), via electromagnetic clutch 32, to rotate drive shaft 11. Rotor 16 rotates synchronously with the rotation of drive shaft 11. Inclined plate 20 rotates synchronously with the rotation of rotor 16. A pair of sliding shoes 22 slide on the surfaces of the radial outer portion of rotated inclined plate 20, while moving reciprocally in the front/rear direction. Piston rod 23, which holds sliding shoes 22, and piston 24, which is formed on the rear end portion of piston rod 23, also move reciprocally in the front/rear direction in cylinder bore 7. By the reciprocal movement of each piston 23, fluid introduced into suction chamber 27 from inlet port 29 is drawn into cylinder bore 7 through suction port 25. This fluid then is compressed in cylinder bore 7 and discharged into discharge chamber 28 through discharge port 26, and then is discharged to an external fluid circuit (not shown) through outlet port 30.\nIn variable displacement inclined plate compressor 1xe2x80x2, when the thermal load of the external fluid circuit increases, the pressures in suction chamber 27 and communication path 9 increase, bellows 31 shrinks, and crank chamber 10 communicates with communication path 9. Blowby gas, that has leaked from cylinder bore 7 into crank chamber 10 through the sliding portion between piston 24 and cylinder bore 7, is released into suction chamber 27 through communication path 9. As a result, the pressure in crank chamber 10 becomes nearly equal to the pressure in suction chamber 27.\nReferring to FIG. 2, in the compression process, moment Ml, which causes inclined plate 20 to rotate in the clockwise direction around pin 21, is generated by the compression reactive force applied to piston 24. Oppositely, moment M2, which causes inclined plate 20 to rotate in the counterclockwise direction around pin 21, is generated by the expanding force of coil spring 19. Therefore, when the thermal load of the external fluid circuit and the pressure in suction chamber 27 increase, the compression reactive force increases to a condition of M1 greater than M2. Consequently, inclined plate 20 rotates clockwise, the stroke of piston 24 increases, and the displacement of variable displacement inclined plate compressor 1xe2x80x2 increases.\nWhen the thermal load of the external fluid circuit decreases, the pressures in suction chamber 27 and communication path 9 decrease, bellows 31 expands, and the communication between crank chamber 10 and communication path 9 is interrupted. The pressure in crank chamber 10 becomes higher than the pressure in suction chamber 27 due to blowby gas introduced into crank chamber 10. Consequently, during the compression process, moment M3, which causes inclined plate 20 to rotate in the counter-clockwise direction around pin 21, is generated by the pressure of blowby gas in crank chamber 10. When the thermal load of the external fluid circuit and the pressure in suction chamber 27 decrease, the compression reactive force decreases, such that M1 less than M2+M3. Consequently, inclined plate 20 rotates counter-clockwise, the stroke of piston 24 decreases, and the displacement of variable displacement inclined plate compressor 1xe2x80x2 decreases.\nIn known variable displacement inclined plate compressor 1xe2x80x2, when compressor 1xe2x80x2 is operated at a low thermal load condition, (a) a resultant axial force, due to the pressure of blowby gas in crank chamber 10, is applied to drive shaft 11 via piston 24, and (b) a resultant force, due to the compression reactive force, is applied to drive shaft 11 via piston 24, balance as time averaged values. In such a condition, the contact/release between rotor 16 and thrust bearing 17 is repeated and drive shaft 11 vibrates, in its axial direction, synchronously with the periodic increase/decrease of the compression reactive force caused by the repeated suction/compression processes. Thus, the contact/release between rotor 16 and thrust bearing 17 and the axial vibration of drive shaft 11 produce noise.\nAccordingly, it is an object of the present invention to provide an improved structure for a variable displacement compressor that prevents an axial vibration of a drive shaft even when operated at a low thermal load condition.\nTo achieve the foregoing and other objects, a variable displacement compressor according to the present invention is provided. The variable displacement compressor comprises a cylinder block having a plurality of cylinder bores and a center bore, a drive shaft extending through a crank chamber, wherein one end portion of the drive shaft is inserted into the center bore and is rotatably supported in the center bore via a radial bearing. An adjusting screw is threaded into the center bore to axially support one end of the drive shaft, via a thrust bearing, and to adjust an axial gap of the thrust bearing. An inclined plate is provided around the drive shaft and rotated synchronously with the drive shaft at an inclined angle, variably controlled relative to an axis of the drive shaft. A plurality of pistons are provided in the cylinder bores and moved reciprocally with the rotation of the inclined plate. A cylinder head is provided on an axial end of the cylinder block via a valve plate and has therein a suction chamber and a discharge chamber for fluid. Compressor displacement is controlled by the inclined angle of the inclined plate, or the pressure in the crank chamber, which thereby controls the stroke of the pistons. The compressor comprises a belleville spring interposed between the thrust bearing and the adjusting screw. recessed portion may be defined on the adjusting screw, and the belleville spring may be in the recessed portion.\nIn the variable displacement compressor, threading the adjusting screw, adjusts the axial gap of the thrust bearing and compresses the belleville spring, which applies a reactive force of the compressed belleville spring to the drive shaft via the thrust bearing. The reactive force of the belleville spring acts in the same direction as that applied with the compression reactive force in the compressor, relative to the drive shaft. Therefore, if the elastic modulus of the belleville spring and the threading degree of the adjusting screw are set at adequate values, and the reactive force of the belleville spring is set at a proper value, even when the compressor is operated at a low thermal load condition, the sum of the time averaged resultant force of the compression reactive force and the reactive force of the belleville spring, that is applied to the drive shaft, is greater than the resultant force of the axial force due to blowby gas in the crank chamber that is applied to the drive shaft. Consequently, the vibration of the drive shaft, that is generated synchronously with the periodic increase/decrease of the compression reactive force due to the repeated suction/compression processes, is prevented.\nIn the present invention, a belleville spring is used to obtain such an advantage. Increase of the axial dimension of the compressor is suppressed by interposing the belleville spring between the thrust bearing and the adjusting screw, as compared to interposing a coil spring therebetween.\nIn a preferred embodiment of the present invention, a recessed portion is defined on the adjusting screw, and the belleville spring is contained in the recessed portion. Therefore, increase of the axial dimension of the compressor is further suppressed."}
{"text": "This invention relates to an apparatus for placing horizontally oriented plates into stacks at a high rate of speed and in a delicate manner.\nIn the production of storage batteries horizontally oriented plates leave an enveloping apparatus serially or in small groups, and these groups must be placed into stacks and the stacks oriented vertically for further processing. A prior art apparatus for accomplishing this objective is shown in Johnson et al., U.S. Pat. No. 4,824,307. Johnson et al. utilizes a series of paddles on a conveyor that comprises a pair of endless chains. The paddles are arranged in pairs that define pockets which carry the stacks. The stacks are formed when a pair of paddles is placed at the end of the conveyor so that the pocket is open and the plates can be dropped into it off the end of a conveyor that carries them horizontally. When a stack is formed the chains are advanced to orient the newly-formed stack vertically and position the next pair of paddles in the proper position to receive plates. This apparatus has several shortcomings which limits its ability to handle delicate plates and to stack plates at high speeds.\nAs the plates are dropped off of the end of the conveyor they rotate to near the desired vertical orientation and this rotation requires a long drop which creates an impact upon landing that can damage delicate plates. In addition, the long drop causes the plates to float which results in their being misaligned when the stacks are formed. Thus, the device disclosed in Johnson et al. has an alignment apparatus which adds to its overall expense.\nIn addition, after each stack is formed the paddles advance, with the newly-formed stack being rotated to a completely vertical orientation and the next set of paddles is placed at the end of the conveyor to form an opened pocket. This advancement must be accomplished while the area above the paddles is free of plates. Since the advancement takes some time to accomplish there is a limit to how fast the plates can be fed to the stacking apparatus on the conveyor and still not have interference.\nThe subject invention overcomes the foregoing shortcomings and limitations of the prior art by providing a conveyor system for transporting horizontally oriented thin rigid plates, in groups of one or more plates, in a manner such that the groups can be displaced from the conveyor system by placing a stop in front of them so as to prevent their continued movement when the conveyor system continues to operate. Thus, the plates displaced from the conveyor drop into stacks while they remain horizontal, thereby permitting the drop distance to be minimized.\nIn a preferred embodiment of the invention this is accomplished by making the space between the chains in the conveyor system greater than the width of the plates and placing inwardly extending support pads on the chains that engage the side margins of the plates. The support pads are placed in spaced-apart sets with each set having a length approximately equal to the height of the plates, and each set being separated from its adjacent set by a distance that is sufficient to allow the plates to pass between them.\nIn one embodiment, the stop is the paddles of the device disclosed in Johnson et al. In another embodiment the stop is located on the piston of a piston cylinder that is arranged so that the stop engages the group of plates when the piston cylinder is extended and does not engage the group when the piston cylinder is retracted.\nThe invention also includes a sorting apparatus that permits selected groups of plates to be removed from the conveyor system before the groups are displaced to form the stack. Thus, reject plates can be removed while the plates are being processed.\nThe invention also includes provisions for varying the distance between the conveyor system and the device that collects the displaced plates. As a result, the collection device can be lowered by a distance equal to the thickness of a group of plates each time a group is displaced so that every group drops the same distance.\nFinally, the invention provides a transfer conveyor that operates at a lower speed than the conveyor system and receives some or all of the groups of plates and delivers them to the stacking apparatus. This causes the groups of plates to arrive at the stacking apparatus at a slower speed than they would if they came directly from the conveyor system. Therefore, there is more time for the completed stacks to be moved out of the way and not interfere with the arriving groups of plates of the next stack, and the apparatus can be operated at a higher rate of speed.\nAccordingly, it is a principal object of the subject invention to provide a plate stacking apparatus in which plates transported in spaced-apart, horizontally oriented groups on a conveyor are slid off the conveyor and collected in stacks while remaining horizontal, thereby only requiring the plates to drop a short distance.\nIt is a further object of the present invention to provide such a device in which the chains of the conveyor that carry the plates are spaced apart by a distance that would permit the plates to drop between them and has inwardly projecting supports that engage the edges of the plates.\nIt is a further object of the subject invention to provide such a device in which each group of plates in a stack drops the same distance as every other group.\nIt is a still further object of the subject invention to provide such a device in which the plates are presented to the stacking apparatus which orients them vertically at a speed that is slower than the speed at which they are transported on the conveyor.\nIt is a still further object of the subject invention to provide such a device which permits defective groups of plates to be removed from the conveyor before they reach the stack so that a full size defect-free stack is always formed."}
{"text": "Most food safety recommendations are for cooking foods to particular temperatures based on the identity of the food item. This requires that the user periodically check the temperature of the food being cooked to see if the desired temperature has been reached, and depending upon the recommendation being used, begin a timer for a period of time while keeping the food at the particular temperature for that time period.\nMost recipes on the other hand are for cooking a food for a particular time period in an environment of a certain temperature on the theory that the food being cooked will reach the needed temperature shortly before or at the expiration of the time period. This approach cause problems, because the oven or other cooking device typically do not obtain the correct cooking temperature or the user does not provide sufficient heat to the food, for example, during stove top cooking or grilling.\nToday given the increase in utility costs, it is desirable to not place a hot food directly into the refrigerator. In addition, however, if a food remains out in or at room temperature for to long or is in the temperature range of 70° F. to 140° F. for over two hours, then a variety of bacteria will start to grow that may cause food borne illnesses for the consumers of the food."}
{"text": "In the manufacture of integrated circuits (ICs), devices are formed on a wafer and connected together by multiple conductive interconnection layers. These conductive interconnection layers are formed by first forming gaps, like trenches and vias in a dielectric layer and then filling gaps with a conductive material.\nThe conductive material is usually formed within the gaps by an electrochemical plating process (ECP process). A barrier layer is firstly formed within the gaps in the dielectric layer. A seed layer is then formed over the barrier layer. The remaining space of the gaps is filled in succession with the conductive material. Then a planarization is performed to remove excess conductive material."}
{"text": "All publications herein are incorporated by reference to the same extent as if each individual publication or patent application was specifically and individually indicated to be incorporated by reference. The following description includes information that may be useful in understanding the present invention. It is not an admission that any of the information provided herein is prior art or relevant to the presently claimed invention, or that any publication specifically or implicitly referenced is prior art.\nMass spectrometry addresses two key questions: (1) “what's in the sample?” and (2) “how much is there?”. Both questions are addressed in the instant application. Several of the embodiments described herein focus on the first question; that is, identification of the components in a mixture. Embodiments of the present invention relate to software that has demonstrated substantial improvements in mass accuracy, sensitivity and mass resolving power. Certain of these gains follow directly from estimation and modeling of ion resonances using a physical model described by Marshall and Comisarow. Other embodiments described herein focus upon applications of estimation and modeling of the phases of ion resonances. Such methods can be divided into functional groups: phase-based methods, calibration, adaptive data-collection strategies, and miscellaneous auxiliary functions.\nThe traditional approach to analysis of Fourier transform mass spectrometry (“FTMS”) spectra is bottom-up. Resonances are detected in the spectra, from which inferences are made about the composition of the analyzed sample. Most of the embodiments described herein involve approaches to bottom-up analysis. Key steps in bottom-up analysis of FTMS data are detection and estimation of ion resonances, mass calibration, and identification. Various embodiments of the present invention involve reducing the 4 MB of data representing an FTMS (MS-1) spectrum to a list of candidate elemental compositions for each detected peak with probabilities assigned to these identities and abundance estimates. The essential information represents a data reduction of roughly three orders of magnitude relative to the unprocessed spectrum. In the bottom-up approach to data analysis, peaks are detected and characterized by estimation first, and then knowledge about the sample is used to calibrate and identify the components. The ability to perform these calculations in real-time creates exciting possibilities for adaptive workflows that actively direct acquisition of optimally informative data."}
{"text": "This invention relates to modulation techniques for transmitting information over fiber optic links.\nIn digital communications systems, bandwidth efficient modulation techniques are used to maximize the amount of information, in the form of digital voice, video and data channels, transmitted across a medium such as a fiber optic link. They do so by reducing the spectral bandwidth required for the transmission of each such digital data channel.\nSome bandwidth efficient modulation techniques involve encoding two or more data bits in each information symbol transmitted across the digital communications link, instead of the familiar single bit binary code. The number of bits encoded by each symbol, and hence the number of possible bit patterns, determines the number of transmission states, or levels, possible for the symbol to take. A four-level code, where each transmitted symbol represents one of four possible symbol states, effectively encodes each of the four possible patterns for a pair of data bits in each transmitted symbol.\nOne technique for transmitting a pair of data bits with a four-level code is quadrature phase shift keying (QPSK) modulation, where each pair of bits is encoded during each symbol period as one of four possible phases of a transmitted carrier signal. Where each pair of bits is derived sequentially from the same data stream, the bit period for the data stream is one half the symbol period of the resulting QPSK signal. However, where each pair of bits is derived from two parallel data streams, the bit period for each data stream is equal to the symbol period of the resulting QPSK signal.\nA technique related to QPSK modulation is differential QPSK (DQPSK) modulation, where each pair of data bits is encoded as one of four possible changes of phase of a transmitted carrier signal. Often, a differential QPSK demodulator for recovering the data bits from the transmitted carrier signal requires circuitry to perform an inverse tangent operation (tan.sup.-1), a data bit synchronization operation, and a four-level logical decision, all operating at data rates up to and beyond 200 Mb/sec.\nIn optical fiber communications systems, optical fiber transmission links typically rely on an optical modulator to modulate a laser light source with the information to be transmitted across the fiber. One type of optical modulator phase modulates the laser light source with one or more microwave subcarrier signals, each itself modulated with digital information. An optical heterodyne detector and subcarrier demodulator at the receiver end of the optical link detects the digital information from each of the modulated subcarrier channels."}
{"text": "The present invention relates to a system and method for adjusting the pressure in an inflatable object. More particularly, the present invention relates to a system and method for adjusting the pressure in an air bed in less time and with greater accuracy.\nAdvances made in the quality of air beds having air chambers as support bases have resulted in vastly increased popularity and sales of such air beds. These air beds are advantageous in that they have an electronic control panel which allows a user to select a desired inflation setting for optimal comfort and to change the inflation setting at any time, thereby providing changes in the firmness of the bed.\nAir bed systems, such as the one described in U.S. Pat. No. 5,904,172 which is incorporated herein by reference in its entirety, generally allow a user to select a desired pressure for each air chamber within the mattress. Upon selecting the desired pressure, a signal is sent to a pump and valve assembly in order to inflate or deflate the air bladders as necessary in order to achieve approximately the desired pressure within the air bladders.\nIn one embodiment of an air bed system, there are two separate air hoses coupled to each of the air bladders. A first air hose extends between the interior of the air bladder and the valve assembly associated with the pump. This first air hose fluidly couples the pump to the air bladder, and is structured to allow air to be added or removed from the air bladder. A second hose extends from the air bladder to a pressure transducer, which continuously monitors the pressure within the air bladder. Thus, as air is being added or removed from the air bladder, the pressure transducer coupled to the second hose is able to continuously check the actual air bladder pressure, which may then be compared to the desired air pressure in order to determine when the desired air pressure within the bladder has been reached.\nIn another embodiment of an air bed system, there is only a single hose coupled to each of the air bladders. In particular, the hose extends between the interior of the air bladder and the valve assembly associated with the pump, and is structured to allow air to be added or removed from the air bladder. Instead of having a second hose with a pressure transducer coupled thereto for continuously reading the pressure within the air bladder, a pressure transducer is positioned within a chamber of the valve assembly. Once the user selects the desired air pressure within the air bladder, the pressure transducer first senses a pressure in the chamber, which it equates to an actual pressure in the air bladder. Then, air is added or removed from the bladder as necessary based upon feedback from the sensed pressure. After a first iteration of sensing the pressure and adding or removing air, the pump turns off and the pressure within the chamber is once again sensed by the pressure transducer and compared to the desired air pressure. The process of adding or removing air, turning off the pump, and sensing pressure within the chamber is repeated for several more iterations until the pressure sensed within the chamber is within an acceptable range close to the desired pressure. As one skilled in the art will appreciate, numerous iterations of inflating and deflating the air bladder may be required until the sensed chamber pressure falls within the acceptable range of the desired pressure.\nThus, while this second embodiment of an air bed system may be desired because it minimizes the necessary number of hoses, it is rather inefficient in that numerous iterations may be required before the sensed pressure reaches the desired pressure. Furthermore, the pump must be turned off each time the pressure transducer takes a pressure measurement, which increases the amount of time that the user must wait until the air bladder reaches the desired pressure.\nTherefore, there is a need for an improved pressure adjustment system and method for an air bed that is able to minimize the amount of time and the number of adjustment iterations necessary to achieve a desired pressure in an air bladder, while also increasing the accuracy of the actual bladder pressure."}
{"text": "Smart phones and electronic tablets provide a convenient way to interact with businesses, especially if those businesses provide an application that allows customers to conduct transactions with the businesses via the applications. Thus, what is desired is a system and method for providing customer service via the applications in a safe and convenient manner."}
{"text": "Systems for presenting a list of items, such as email messages, are generally well known. Typically, only a small portion of text associated with these items is displayed in the list itself. If a user wishes to view more information, the user must select one of the items to additional information would be displayed in a separate area.\nOften times, messages may be related, but appear scattered as separated items, making it difficult for a user to follow a complete discussion. Furthermore, in typical systems, a user wishing to mark a message for a particular action, such as urgent or delete, must move to a different area of the screen in order to accomplish this.\nThere are many situations in which people would like to view extended amounts of text, but are limited by the amount of display area or “screen real estate” available to them. For example, a user looking at a list of items might want to see or preview the detailed content text of one of the items, without looking to or covering up other areas of the display.\nAnother problem is that a user wishing to follow a discussion must sort through a list of scattered messages. This can make it difficult for users to locate related items, to notice if an issue raised in a message has already been handled by someone else, or to follow a discussion at all. This can be particularly difficult when using a portable device with a small display as there is less opportunity to see related items in a large displayed list. Furthermore, due to difficulties in entering text on many such devices, the cost of responding to a message that has already been handled by someone else is high.\nAnother problem is that conventional systems require a user wishing to mark an item for later action to move the cursor to a different region of the display or use keyboard controls. This is especially problematic in situations where screen real estate is limited, such as with handheld portable devices.\nOther limitations and problems also exist."}
{"text": "1. Field of Invention\nThe present invention relates to fuel injection valves for use in internal combustion engines.\n2. Description of Related Art\nVarious types of fuel injection valves have been hitherto proposed and put into practical use particularly in the field of automotive internal combustion engines. One of them is disclosed in Laid-open Japanese Patent Application (Tokkai) 2000-8990, which generally comprises a valve element arranged in a casing to move between open and close positions and an electromagnetic coil arranged to actuate the valve element to move between the open and close positions in accordance with energization/de-energization of the electromagnetic coil. Upon opening of the valve element, a pressurized fuel in a fuel passage of the casing is injected into a given space, such as a combustion chamber, air intake passage or the like, of the internal combustion engine.\nHowever, due to inherent construction, the fuel injection valves of the type disclosed by the above-mentioned Japanese Application tend to fail to exhibit a satisfied dimensional stability of the valve element. That is, for producing the valve element, a press working is employed to shape a given portion of the valve element, which however tends to induce a deformation or swelling of the given portion to which the pressing force is actually applied. As is known, such deformation or swelling (viz., dimensional poorness) induces a non-smoothed movement of the valve element and thus tends to exhibit a poor responsiveness of the same upon energization/de-energization of the electromagnetic coil. If, for smoothing the movement, a finish machining working is additionally employed for finely shaping the valve element, productivity of the fuel injection valve is lowered and thus cost performance of the same becomes poor."}
{"text": "1. Field of the Invention\nThe present invention relates to novel oligonucleotide analogues, which exhibit antisense or antigene activity having excellent stability, or exhibit excellent activity as a detection agent (probe) for a specific gene or as a primer for starting amplification, and to novel nucleoside analogues which are intermediates for their production.\n2. Background Information\nOligonucleotide analogues, which have excellent antisense or antigene activity and which are stable in the body are expected to be useful pharmaceuticals. In addition, oligonucleotide analogues having a high degree of stable complementary chain formation ability with DNA or mRNA are useful as detection agents for a specific gene or as primers for starting amplification.\nIn contrast, naturally-occurring oligonucleotides are known to be quickly decomposed by various nucleases present in the blood and cells. In some cases, naturally-occurring oligonucleotides may not have sufficient sensitivity for use as detection agents for specific genes or as primers for starting amplification due to limitations on their affinity with complementary base sequences.\nIn order to overcome these shortcomings, various non-naturally-occurring oligonucleotide analogues have been produced, and have been attempted to be developed for use as pharmaceuticals or detection agents for specific genes. Namely, known examples of such non-naturally-occurring oligonucleotide analogues include those in which an oxygen atom attached to a phosphorus atom in a phosphodiester bond of an oligonucleotide is replaced with a sulfur atom, those in which said oxygen atom is replaced with a methyl group, those in which said oxygen atom is replaced with a boron atom, and those in which a sugar moiety or base moiety of an oligonucleotide is chemically modified. For example, ISIS Corp. has developed thioate-type oligonucleotide ISIS2922 (Vitravene) as a therapeutic agent for human cytomegalovirus retinitis and ISIS2922 has been put on the open market in the United States.\nHowever, in consideration of the potency of the antisense or antigene activity in the above non-naturally-occurring oligonucleotide analogues, namely the ability to form a stable complementary chain with DNA or mRNA, stability with respect to various nucleases, and the manifestation of adverse side effects due to non-specific bonding with various proteins in the body, there has been a need for a non-naturally-occurring oligonucleotide analogue having even better stability in the body, a low incidence of adverse side effects and a high ability to form complementary chains."}
{"text": "1. Field of the Invention\nThe present invention relates in general to a cathode ray tube, more particularly, to a cathode ray tube mounted with a reinforcing band for improving landing stability of electron beams and howling of shadow mask and reinforcing an explosion-proof function.\n2. Discussion of the Background Art\nFIG. 1 illustrates the structure of a related art cathode ray tube, and FIG. 2 is a cross-sectional view taken along line A–A′ in FIG. 1.\nAs shown in FIGS. 1 and 2, the related cathode ray tube includes a panel 3 having a fluorescent screen 6 formed on an inner surface, a funnel 5 connected to a seal edge portion 4 of a skirt portion 2 of the panel 3, an electron gun 14 mounted on a neck portion 12 of the funnel 5, emitting electron beams 13, deflection yokes 15 for deflecting the electron beams 13 emitted from the electron gun 14, a shadow mask 8 for selecting colors of the electron beams 13, and a mask frame 9 for supporting the shadow mask 8.\nThe mask frame 9 is connected to a stud pin 10 formed on the skirt portion 2 of the panel 3 by means of mask springs 11.\nAlso, a reinforcing band 17 straps around an outer peripheral portion of the skirt portion 2 of the panel 3 in order to prevent CRT implosion by external impacts, and a lug 18 to be incorporated in a TV set is also formed on the outer peripheral portion of the skirt portion 2 of the panel.\nThe Z—Z line in FIG. 1 indicates a tube axis.\nAs depicted in FIG. 3, the reinforcing band 17 is usually mounted on a portion where the stud pins 10 are formed.\nHowever, when the reinforcing band 17 is mounted on where the stud pin 10 is formed, a howling phenomenon occurs to the shadow mask 8 due to voice vibrations generated from a speaker mounted on the outside of the cathode ray tube.\nTo be more specific, the voice vibrations generated from the speaker are transmitted to the reinforcing band 17, and to the stud pin 10 formed inside of the skirt portion 2 of the panel 3, and eventually to the mask frame 9 and the shadow mask 8 connected to the stud pin 10.\nAs a result, the howling phenomenon occurs to the shadow mask 8, and electron beams do not land stably, eventually deteriorating a picture quality.\nAs FIG. 4 shows, clamping power of the reinforcing band 17 is influenced by the outer peripheral length of the skirt portion 2 of the panel 3 and the inner peripheral length of the reinforcing band 17. The clamping power of the reinforcing band 17 influences in turn picture quality characteristics of the cathode ray tube.\nFIG. 4 illustrates clamping power errors of the reinforcing band. More details on the drawing are now followed.\nReferring to FIG. 4, Lm denotes an inner peripheral length of the reinforcing band 17, ±β denotes tolerance errors in a manufacturing process for the reinforcing band, Lg denotes an outer peripheral length of the skirt portion 2, and ±α denotes tolerance errors in a manufacturing process for the skirt portion. Here, the reinforcing band 17 has a maximum clamping power at a point Pa, and a maximum clamping power at a point Pb.\nGenerally, the reinforcing band 17 is expanded by heating. The reinforcing band 17 is welded to the skirt portion by stress applied thereon when the reinforcing band is strapped around the outer peripheral portion of the skirt portion 2. And, the reinforcing band 17 becomes clamped by cooling. However, depending on temperature and material characteristics, there can be errors, and these errors eventually change the clamping power of the reinforcing band.\nErrors of the clamping power of the reinforcing band 17 contribute to the difference in compressive stress and the deviations in landing of electron beams.\nAfter all, a picture quality of the cathode ray tube is severely deteriorated. This problem becomes worse, however, as shown in FIG. 3, if the reinforcing band 17 is mounted on an outer surface of the skirt portion where the stud pin 10 is formed.\nIn addition, compared to a traditional cathode ray tube having a convex outer surface, an inner surface of radius of curvature (Ri in FIG. 3) of the panel whose outer surface is substantially flat is increased, and this in turn increases the curvature radius of the shadow mask 8, thereby lessening the strength of the shadow mask 8.\nParticularly, in case of using tint glass having less than 51% of light transmittance, the difference between a thickness at the central portion (Tp) of the panel and a thickness at the peripheral portion (Tw) of the panel causes about 15 fL of brightness reduction to the peripheral portion. Thus, the thickness at the peripheral portion (Tw) should be reduced.\nHowever, if the peripheral portion of the panel 3 gets thinner, the radius of curvature of the inner surface of the panel 3 is increased, and this causes an increase of the radius of curvature of the shadow mask 8 disposed at a predetermined distance (Tmp) from the inner surface of the panel 3. When these happen, the strength of the shadow mask 8 is lessened, being more subject to the influence of external impacts.\nFIG. 5 shows different thicknesses of a panel with and without a tint glass.\nSince the thickness at the peripheral portion of a tint glass (Pt) is relatively thinner than a clear glass (Pc), as shown in FIG. 6, the radius of curvature of the inner surface (Ri) of the tilt glass (Pt) is greater than that of the clear class (Pc), and the strength of the shadow mask 8 is lessened even more.\nTherefore, the howling phenomenon due to voice vibrations of the speaker, and clamping power difference (error) of the reinforcing band 17 play a more negative role in landing of electron beams.\nFIG. 7 describes stress being applied to the inside of a cathode ray tube in vacuum state.\nWhen the inside of the cathode ray tube is in vacuum state, compressive stress (σc) is applied to the central portion of the panel 3, and tensile stress (σt) is applied to the peripheral and skirt portions of the panel 3.\nIf the cathode ray tube in such state is broke down by external impacts, fragments of the panel 3 are first sucked into the cathode ray tube because of atmospheric pressure, and split to the outside later.\nThe above phenomenon is called ‘accidental implosion phenomenon’. Thus, the reinforcing band 17 constricts the skirt portion 2 to minimize splitting of glass fragments at the time of CRT implosion.\nFIG. 8 shows how stress changes according to a reinforcing band mounted on the skirt portion.\nReferring to FIG. 8, if the reinforcing band 17 is mounted on a skirt portion 2 of a panel having a predetermined outer surface radius of curvature (Ro), clamping power (P) of the reinforcing band 17 acts upon the skirt portion 2, thereby lessening compressive stress at the central portion of the panel 3.\nMore specifically, as the compressive stress is lessened by mounting the reinforcing band 17, a variation (hv) of the panel before the reinforcing band 17 is mounted thereon is changed to a variation (hb) of the panel after the reinforcing band 17 is mounted thereon.\nWhat happens here is that the compressive stress is cancelled by the clamping power (P) of the reinforcing band 17. Therefore, even when CRT implosion accidentally occurs due to the lessened compressive stress at the central portion of the panel 3, much less fragments are split outward.\nFIG. 9 illustrates how clamping power of a reinforcing band acts upon a skirt portion of a panel.\nAs depicted in FIG. 9, when the reinforcing band 17 is mounted on the skirt portion 2 of the panel 3 having a predetermined radius of curvature (Ro) on its outer surface, the clamping power (P) of the reinforcing band 17 is divided into a vertical component of power (σn) and a horizontal component of power (σp), and their resultant power (σr) is applied along the curvature radius of the panel 3.\nThus, the compressive stress acting upon the central portion of the panel 3 because of atmospheric pressure is cancelled by the resultant power (σr), and as a result, the compressive stress is lessened, as discussed in FIG. 8.\nHowever, the above operation is not always equally applied to a substantially flat outer surface.\nFIG. 10 describes changes in stress after mounting a reinforcing band on a panel having a substantially flat outer surface, and FIG. 11 illustrates how clamping power of a reinforcing band acts upon a panel having a substantially flat outer surface.\nAs shown in FIG. 10, a variation (hb) of compressive stress acting upon a central portion of the panel 3 where the reinforcing band is mounted on the panel having a substantially flat outer surface is rather increased, in contrast with what is shown in FIG. 8.\nThat is, thanks to the clamping power (P) of the reinforcing band, the compressive stress after mounting the reinforcing band 17 is much improved, compared to before mounting the reinforcing band 17.\nThis phenomenon occurs because the clamping power (P) of the reinforcing band 17 is divided into a vertical component of power (σn) and a horizontal component of power (σp), as shown in FIG. 11, and their resultant power (σr) is applied to the inward of a cathode ray tube because the outer surface of the panel 3 is substantially flat.\nThat is to say, it is hard to expect the same effect, namely improving the explosion-proof characteristic of a cathode ray tube, from mounting the reinforcing band 17 on the panel having a substantially flat outer surface with the one on the panel whose outer surface is curved,"}
{"text": "1. Field\nThe following technology relates to a technology that reduces interference in a macro cell communication system when the macrocell and a small cell use the same radio resources.\n2. Description of Related Art\nRecently, an interest in small cell communication systems, such as a femtocells, picocells, and the like, have increased. Unlike macrocells permanently installed by service providers, the small cells may be freely installed by users. At least one macrocell and at least one small cell may construct a hierarchical cell communication system.\nThe hierarchical cell communication system may establish at least one small cell in a predetermined location, to obtain a high capacity of system resources in the predetermined location. However, interference may occur in the macrocell due to the small cell. The interference may reduce system throughput and cause a loss of data."}
{"text": "The present invention relates generally to the testing of logic designs, and more particularly to a method and system for generating a test case specific to a particular design-under-test (DUT) from a generalized or comprehensive set of bus transactions characterizing a bus architecture.\nA process known as verification comprises applying test cases to a DUT and checking the results against expected results to determine whether the DUT functions correctly. Functional verification of logic designs which interface to standardized buses typically involves the use of a software tool known as a bus functional model (BFM). A BFM emulates a bus protocol to apply a test case in simulation to a DUT modeled by a hardware description language (HDL).\nThe BFMs allow verification engineers to write test cases as bus transactions. The bus transactions may be expressed in bus functional language (BFL) statements which are decoded and executed by the BFMs using a logic simulator. For example, a standardized BFM transaction test case for a 32 bit data slave design that uses read and write transactions with address, data, and control signal information could comprise BFL statements appearing as follows (comments indicated by xe2x80x9c//xe2x80x9d):\nIn the foregoing example, read and write are two different transaction types, while address, byte enable, and data are parameters of the respective transaction types. The parameters were used to specify data to be written to an address, and to subsequently read the address to determine whether the data was correctly written.\nIn order to verify that a DUT completely complies with a standardized bus interface, all possible bus transactions must be performed to ensure that the DUT-to-bus interface operates under all legal bus conditions. This process, also known as regression testing, entails the enumeration of all possible and legal transaction combinations to ensure a complete and correct design.\nTypically, test cases for applying all the possible and legal transactions for a DUT-to-bus interface must be specified manually, as BFL statements as shown in the above example. To code the BFL statements, a verification engineer must usually consult a bus architecture specification manual. As the complexity of logic designs increases, the number of transactions which must be applied to a design to fully test it for bus compliance becomes unmanageable by a manual procedure. The procedure becomes unduly time-consuming, and the likelihood that coding errors occur or that needed test cases are omitted increases.\nIn view of the foregoing, a method and system for efficiently generating test cases is called for.\nIn the method and system according to the present invention, a set of transactions characterizing a bus architecture is provided. A subset of the set corresponding to a specific DUT is selected, and a configuration file for verifying the DUT is generated from the selected subset. The configuration file includes rules specific to the DUT, and may be processed by a generator program to generate or enumerate bus transactions, as determined by the rules, for a test case to be applied to the DUT.\nIn an embodiment, a graphical user interface may be used to create the configuration file. The user interface presents a user with the possible bus transactions of a bus architecture, and allows the user to enter inputs corresponding to a particular DUT. The interface processes the inputs to generate a configuration file syntax used by the generator program. The interface automatically inserts parameter combination specifications and rules into the configuration file for the user.\nThe advantages offered by the foregoing include greater ease and speed of designing test cases for verification of designs, and avoiding the errors inherent in a manual procedure. A designer of a test case does not need to consult a bus architectural specification, since the set of transactions characterizing the bus is presented in the user interface. A test case may be designed by simply entering a series of inputs relating to a specific DUT via the interface, which then automatically produces the configuration file."}
{"text": "A number of applications, including certain clock and data recovery (CDR) architectures, require a set of evenly spaced clock waveforms to sample a received data waveform. In these applications, this set of clock waveforms is typically generated by using a voltage-controlled delay line (VCDL). A voltage-controlled delay line 110, as shown in FIG. 1, is an electrical circuit that is comprised of a plurality of sequentially arranged time delay elements, which is often known as a delay chain, driven by a reference clock signal. Thus, an N-cell VCDL generates N output clocks, where each successive output clock is delayed from the input clock by an additional unit delay. The time delays of each of the individual elements that comprise the chain are substantially equal, and are controlled by an analog voltage or current.\nIt is typically desired that the total delay through the chain of N delay cells be equal to the period, T, of the input reference clock signal, often referred to as the injection clock, INJ. As shown in FIG. 2B, this implies that the time delay between the rising edges (and the delay between the falling edges) of adjacent output clocks is equal to T/N, and also that the rising (and falling) edges of the output of the Nth delay cell, often referred to as the return clock, RET, are aligned to the rising (and falling) edges of the injection clock, INJ. In addition, if the duty cycle of CLKIN is exactly 50%, for any two output clocks that are separated by exactly N/2 cells in the chain, the rising (falling) edge of the earlier clock is aligned to the falling (rising) edge of the later clock.\nThe time delay through an electrical delay cell varies significantly due to variations in the manufacturing process (P) and with the power supply voltage (V) and operating temperature (T), often collectively referred to as PVT variations. Thus, most applications of VCDLs require a continuous feedback loop, known as a delay-locked loop (DLL), that senses the alignment between the edges of relevant clocks and varies the control voltage to increase or decrease the delay of each stage in the chain, as appropriate, to hold the total VCDL delay equal to T. Typically, a phase detector in the DLL senses the alignment of the edges of certain VCDL clocks. For example, the phase detector may sense the alignment between the rising edge of the injection clock and the rising edge of the output of the final delay cell in the chain (i.e., the return clock).\nDuty cycle distortion is a source of impairment in a received signal. The duty cycle of a phase locked loop (PLL) is the percentage of time that the output clock signal has a given value. A PLL should typically demonstrate a 50% duty cycle, such that the output clock signal should alternate between two amplitude values, each for 50% of the total duration. Duty cycle distortion (DCD) arises due to device mismatches and due to variations of the differential signal paths for clock and data. The target 50% duty cycle feature is particularly important for high-speed applications where both positive and negative edges are used. In addition, each delay element in a VCDL chain acts as a low pass filter. Thus, the corresponding frequency-dependent attenuation will cause the DCD to get progressively worse with each delay element, often referred to as DCD build up.\nA need exists for an improved control mechanism for starting up a VCDL with a wide capture range. A further need exists for a VCDL edge alignment process that is substantially immune from PVT variability."}
{"text": "Automobiles and trucks are commonly fitted with towing hitches to allow them to pull trailers of assorted types. The trailers may support boats, off road vehicles, or campers, or may be for transport of industrial or home equipment. In any event, the value of the trailer and the towed goods can be considerable, presenting an attractive target for thieves. However, because of the large size of the automobile-trailer assembly, it must at times be parked outside and will occasionally be left unattended. Automobiles are provided with numerous well known impediments to theft, i.e.: door locks, ignition locks, alarms, etc. To prevent the trailer from being separated from the hitch, the coupler arm of the trailer is typically locked to the hitch ball which is supported on the vehicle towing hitch.\nThe hitch ball may be attached directly to an attaching member of the pulling vehicle which may be the bumper of the pulling vehicle or a special extending tongue that is affixed to the bumper or underbody of the vehicle. The hitch ball has a spherical upper section with a downwardly protruding threaded shaft which extends through the vehicle hitch. Because different trailers require hitch balls of different diameter, for example from 1-7/8 to 2-5/8, the hitch ball must be separable from the vehicle hitch. The shaft of the hitch ball is held in place by a nut threaded beneath the tongue or the bumper of the vehicle. The towed vehicle trailer arm has a cup or socket which fits over the hitch ball and allows a pivoting and rotating connection between the trailer and the vehicle. The trailer arm will typically have structure which allows it to be secured and locked to the hitch ball.\nWhile trailer hitch socket couplers are commonly permanently attached to the trailer, hitch ball couplers are specifically designed so that a hitch ball can be quickly and easily removed and replaced by a hitch-ball of the correct size to properly fit a trailer hitch socket coupler with a different sized socket. Thieves, too, often readily exploit these hitch ball coupler design features and in a matter of minutes can remove the hitch ball still clamped and locked into place to the trailer and reposition the hitch ball shaft in a hitch on a vehicle under their control, making it possible to rapidly drive off with the trailer and its contents.\nTo prevent the hitch ball from being separated from the towing vehicle hitch, it has heretofore been necessary to in some way lock the hitch ball independently, for example by a padlock inserted through the shaft beneath the hitch tongue. Not only does an additional lock add to the cost of the towing assembly, it represents an inconvenience to the driver, who must carry an additional key and take the additional time required to unlock the hitch ball lock each time the hitch ball is separated from the towing vehicle hitch.\nAn unlocked hitch ball not only accommodates thieves, it also presents the possibility that the securing nut may work its way off the shaft, leading to accidental disconnection while travelling. Oftentimes, additional coupling devices such as chains or secondary clamping devices and locks are employed to prevent accidental disconnection.\nA trailer which is parked without a hitch ball received within the trailer arm is also susceptible to being mounted on a hitch ball or other protrusion from a towing vehicle and being removed by thieves. To prevent the trailer and from being used in this manner, hitch balls without protruding shanks, such as those sold under the name TRAILER COP have been employed. Such a device fills the cavity in a trailer arm and prevents insertion of any supporting element. However, the hitch ball without a shaft must be removed and replaced with a hitch ball with a shaft when it is desired to connect the trailer to a towing vehicle.\nWhat is needed is a hitch ball that is relatively simple, inexpensive, easy to manufacture, and which is securable to the towing vehicle and trailer without the requirement of a separate lock."}
{"text": "1. Technical Field\nThe present invention relates to networks of devices that can be connected using wireless links. Specifically, embodiments of the present invention pertain to a method and system for selecting and connecting to wireless access point within a local area network.\n2. Background Art\nPrior art FIG. 1 depicts a radio frequency (RF) wireless microphone 11 that is sending audio, via a RF communication link 100, to a wireless audio access point 18. This application is suitable for a teacher, using this wireless RF microphone 10, to address students in a classroom setting. Audio distribution using wireless RF communication is well understood in the prior art. Similar transmission methods using known RF communication protocols are also well understood in the art. For example, Bluetooth® is one such method. Bluetooth is a trademark of Bluetooth SIG, Inc. Wi-Fi® devices communicate using the IEEE 802.11 standards. Wi-Fi is a trademark of the Wi-Fi Alliance. It is also well known in the prior art to use IEEE 802.11 in conjunction with Internet Protocol to wirelessly transmit audio data, such as in Voice over Internet Protocol (VoIP) applications.\nOne recent advance in wireless RF communication is Wi-Fi Direct, a wireless standard that allows Wi-Fi devices to communicate directly with each other without the need for a router. This peer-to-peer communication standard can be implemented in any Wi-Fi device. Those skilled in the art will recognize the advantages of transmitting audio data directly from a wireless microphone to a wireless audio access point in accordance with the Wi-Fi Direct specification.\nWireless RF communication can be established by a device, such as a RF wireless microphone 11, by first broadcasting an inquiry message 101 to a wireless audio access point 18 which responds by broadcasting an inquiry response message 102.\nFIG. 2 depicts a portion of a school building having several classrooms with the range of a typical wireless device thereupon superimposed. This clearly shows that in a school setting, the range 210 of a typical wireless device, such as wireless audio access point 18, located in a first classroom 21 extends well beyond the confines of that classroom. As shown, the wireless audio access point 18 located in the first classroom 21 would try to establish RF communication 100, with any compatible wireless devices located in a second classroom 22, a fourth classroom 24, and a fifth classroom 25. The wireless audio access point 18 also potentially interacts with some portion of the devices in a third classroom 23, a sixth classroom 26, a seventh classroom 27, and an eighth classroom 28. In fact, in the typical school layout illustrated only devices in a ninth classroom 29 would not interact with the wireless audio access point 18 in the first classroom 21.\nThe prior art does not provide an effective means of limiting wireless communications to the confines of a single classroom. The problems associated with employing a wireless microphone device in a multiple classroom setting are compounded when multiple wireless microphones are employed. For example, envision a classroom setting with both a teacher microphone and multiple student microphones. The solution to the original problem must now take into account the additional issues of interference, coordination and interoperability."}
{"text": "1. Field of the Invention\nThe present disclosure relates to image processing, and more particularly to a system and method for detecting an object in an image while eliminating false positive detections.\n2. Description of Related Art\nThe pulmonary embolism (PE) is defined as a thrombus, which is generally recognized as dark regions within enhanced pulmonary arteries in computed tomography angiography (CTA) images. For PE detection, it is difficult to distinguish PE from various PE look-alikes (false positives) including flow voids in veins and arteries, lymphoid tissues, streak artifacts near superior vena cava, partial volume artifacts at vascular bifurcations, etc.\nTherefore, a need exists for a system and method for object detection including a feature based approach for automatically removing the false positives, while preserving the true positives."}
{"text": "1. Field of the Invention\nThe present invention is related to power supply regulators for radio frequency applications and more particularly to a shunt regulator for high frequency applications.\n2. Background Description\nVoltage regulators are well known. An ideal voltage regulator provides a constant voltage regardless of load. Thus, the voltage regulator provides the same voltage under no load (at no current) as it does fully loaded. Current used by circuits in the complementary insulated gate field effect transistor (FET) technology, commonly known as CMOS, primarily, is switching current with negligible static (or DC) current flow.\nCMOS circuit current flow usually occurs only during switching, primarily, either to charge or discharge the circuit\"\"s load (capacitance). Thus, digital circuits that are synchronized by a common clock signal may exhibit sporadic episodes of very high switching current, e.g., from a counter switching from FFFF16 to OOOO16. By contrast, typical radio frequency circuits, such as may be used in a radio telephone or cellular phone, exhibit relatively uniform switching and, therefore, have relatively uniform current. Variations between no load and full load current, such as may occur with digital circuits, can cause large switching noise that must be filtered to prevent errors in CMOS analog circuits on the same integrated circuit substrate or chip. Shunt regulators may be used to reduce switching and other current related noise.\nA typical shunt regulator includes an alternate current path for regulator current, the regulator supplying constant current during no load conditions as well as during full load. The parallel current path maintains an effective load such that even as the load varies, the regulator supplies constant (full load) current with excess current being shunted through the parallel path. Many state of the art shunt regulators are designed to maintain effective constant current. That unused portion of the full load current in excess of the load current is shunted through a shunt device in the shunt regulator and so, wasted.\nHigh frequency circuits, and especially radio frequency (RF) circuits, are very sensitive to noise. Switching noise from digital circuits can easily couple into radio frequency circuits thereby, degrading circuit performance. As higher degrees of integration are being achieved, larger numbers of complex high frequency circuits are being integrated onto a single integrated circuit chip. Further, as digital circuit performance improves, digital functions are also being combined with RF functions onto monolithic integrated circuit chips integrating more and more digital and RF circuits onto the same chip and resulting in an increased number of local potential noise sources on a given chip. To mitigate this problem on-chip voltage regulators may be used to provide separate isolated voltage supplies for digital and for RF circuits, thereby, isolating digital switching currents from input power supplies and from other on-chip RF circuits. This approach significantly reduces the digital switching noise that might otherwise couple into the RF circuits. Typically, shunt regulators are used for this type of isolation.\nA typical prior art shunt regulator is described in U.S. Pat. No. 4,366,432 entitled xe2x80x9cHighly Stable Constant-Voltage Power Source Devicexe2x80x9d to Noro. Noro describes a shunt regulator that is biased to deliver a constant current to a parallel combination of a load and shunt device. This constant current value must be set high enough to supply the maximum load current and any lesser variations thereof with the excess passing through the shunt device. Consequently, Noro\"\"s shunt regulator constantly provides the maximum current and, when there is no load current, i.e., it is unloaded, all of the supply current is shunted through the shunt device. Accordingly, when the load is less than full, Noro wastes some of the power supplied.\nU.S. Pat. No. 5,260,644 entitled xe2x80x9cSelf-Adjusting Shunt Regulator and Methodxe2x80x9d to Curtis describes a shunt regulator which attempts to isolate power supply and load, automatically adjusting supply current to compensate for load variations. Unfortunately, the Curtis shunt regulator also consume and excess power when it is not fully loaded. The Curtis shunt regulator varies shunt current to compensate for changes in load currents such that fluctuation of total supply current are only a small fraction of fluctuations of load current.\nThese prior art shunt regulators are used to provide isolation between a chip power supply and a digital circuit portion of the chip load. These shunt regulators can tolerate rapidly varying load current. Unfortunately, these typical prior art shunt regulators suffer from excessive power consumption, while only providing limited isolation between digital and RF circuits that are integrated onto the same integrated circuit chip.\nThus, there is a need for improved shunt regulators that more effectively isolate digital circuits from RF circuits integrated onto the same chip, but with a minimum increase in load current."}
{"text": "Integrated circuit technology has revolutionized various fields, including computers, control systems, telecommunications, and imaging. One field in which integrated circuitry is widely used is video imaging. Different types of semiconductor imagers include: charge coupled devices, photodiode arrays, charge injection devices, and hybrid focal plane arrays. Many of these devices include pixels that are arranged in sensor arrays to convert light images into electrical signals.\nExamples of MOS imaging devices are detailed in “A ¼ Inch Format 250K Pixel Amplified MOS Image Sensor Using CMOS Process” by Kawashima et al., IEDM 93–575 (1993), and “A Low Noise Line-Amplified MOS Imaging Devices” by Ozaki et al., IEEE Transactions on Electron Devices, Vol. 38, No. 5, May 1991. In addition, U.S. Pat. No. 5,345,266 to Denyer, titled “Matrix Array Image Sensor Chip,” describes a MOS image sensor. The devices disclosed in these publications provide a general design approach to MOS imaging devices. In addition, MOS approaches to color imaging devices are described in “Color Filters and Processing Alternatives for One-Chip Cameras,” by Paruiski, IEEE Transactions on Electron Devices, Vol. ED-32, No. 8, August 1985, and “Single-Chip Color Cameras With Reduced Aliasing” by Imaide et al., Journal of Imaging Technology, Vol. 12, No. 5, October 1986, pp. 258–260.\nImage sensor circuitry generally includes circuits for performing black level calibration and automatic gain control. Black level calibration attempts to eliminate the portion of the image signal that exists when no light is being detected, thus allowing for a truer zero reference signal so that the later signal processing is improved. After the black level calibration has been performed, automatic gain control amplifies the video signal at a controlled level so as to utilize more of the available signal amplification range. The combination of the black level calibration and automatic gain control can be said to form one type of an “auto black expansion” method. One prior art circuit that addresses issues related to these processes is shown in U.S. Pat. No. 4,187,519 to Vitols et al. FIGS. 1, 2, 3, 4, 5, 6 and 7 of Vitols et al. have been reproduced herein as FIGS. 1, 2A, 2B, 2C, 2D, 2E, and 2F, respectively.\nVitols et al. disclose a circuit that is designed to provide a video contrast expansion system. As illustrated in FIG. 1, a sequence of video pixel data representative of an image is developed by a signal source 11, which may be, for example, a radar, an IR (infrared) sensor or a TV (television) camera. This pixel data is applied to an I-getter (intensity getter) 13 and to a delay circuit 19. An initialize command signal from the source 11 initializes the operation of the I-getter 13 before or during the time that the sequence of pixels is being generated by the source 11.\nThe I-getter 13 searches the input sequence of pixels (designated as new data or ND) to determine bias and gain intensity parameters. A bias function generator 15 is automatically adjusted by the bias parameters to develop pixel bias correction signals. In a similar manner a gain function generator 17, which is similar in structure and operation to the generator 15, is automatically adjusted by the gain parameters to develop pixel correction signals.\nIn its operation the I-getter 13 reduces and transforms the input sequence of pixels into a reduced number of bias parameters and a reduced number of gain parameters. Each of the function generators 15 and 17 smoothly fills in or shades in its associated widely spaced apart parameters (bias or gain) by double linear interpolations to produce its associated correction signals (bias or gain). Each sequence of correction signals is equal to the number of pixels in the input sequence of pixels.\nIt takes time for each of the function generators 15 and 17 to get sufficient parameter data in before it can start computing its associated correction signals. Furthermore, it takes additional time before each of the function generators 15 and 17 compute its associated correction signals. This combined delay time is offset by the delay circuit 19 which synchronizes the time of occurrence of the sequence of pixels at the output of the delay circuit 19 with the time of generation of the pixel bias correction signals and pixel gain correction signals.\nThe computed pixel bias correction signals are respectively subtracted in a combiner or subtractor 21 from the delayed or synchronized sequence of pixels to selectively lower the minimum values in associated groups of pixels in the delayed sequence of pixels to very close to zero. The minimized output of the combiner 21 is respectively multiplied in a multiplier or AGC (automatic gain control) circuit 23 to selectively expand the maximum amplitudes in the associated groups of pixels at the output of the combiner 21 to near the saturation level of the electronics of the system of FIG. 1. The output of the AGC circuit 23 can be applied to a display generator (not shown) to generate an enhanced picture of the image in which the video contrast of the image has been substantially expanded in two directions (minimum and maximum values of contrast).\nThe pixel data at the output of the signal source 11 may have a very narrow dynamic range of contrast, making it extremely difficult to discern objects in areas of lower contrast if directly viewed at the point. However, by the operation just explained, the pixel data at the output of the AGC circuit 23 has its dynamic range of contrast selectively expanded or stretched from near minimum to near maximum, making objects originally in low areas of contrast now clearly discernible when displayed.\nFIGS. 2A–2F illustrate some of the possible operations that the circuit of FIG. 1 can perform on a one-dimensional signal. Such operations are analogous to those performed on a two-dimensional signal since a two-dimensional signal is essentially comprised of a vertical plurality of horizontal one-dimensional signals.\nFIG. 2A illustrates an exemplary one-dimensional signal I(x) comprised of a sequence of 256 pixels. Although I(x) is shown as an analog signal it is a sequence of 256 intensity (I) points of varying amplitudes. Note that I(x) varies over a fairly wide intensity range from an amplitude of 630 (Isat or the intensity saturation level of the system) to an amplitude of about 50 at the low end.\nIt may be desired in some applications to look for small, rather than large, amplitude deviations from nominal. For example, if the intensity variations in the incoming signals I(x) of FIG. 2 represent pixels or intensity points of a scene, very high intensity portions of I(x) would ultimately develop very bright portions in a picture, and very low intensity portions of I(x) would ultimately develop very dark portions of the picture. However, a human observer or operator may not necessarily be interested in very high or very low intensity levels in the incoming signal I(x). Rather, an observer may only be interested in small variations of intensity in I(x), whether contained in very high, very low and/or intermediate levels of intensity in I(x). This is due to the fact that the small variations of intensity in I(x) can define the details of a scene or image sufficiently to possibly enable an observer to identify what is happening or contained in the picture.\nBy using a conventional contrast enhancement technique, an observer may be unable to discern what is contained in picture areas of small intensity variations. It is to the correction of this problem that the circuit of FIG. 1 is directed. To illustrate, assume that the low intensity region of pixels between 100 and 200 in FIG. 2 represents the signal that is coming from the inside of a dark cave. The brightest pixel in the 100–200 region of pixels is dark in comparison with the other high intensity or bright pixels in FIG. 2. An observer may be looking for the glint of light on a gun barrel inside that cave. If the dynamics of the signal in that 100–200 pixel region remain unchanged, an observer would not be able to see that glint of light on that gun barrel. However, if the gain of the signal in the 100–200 pixel region (where the amplitude of the signal is very low) could be increased so that variations of the signal show up to substantially the fullest extent possible without saturating the electronics of FIG. 1, the human observer could more readily determine what was inside the cave. It should be noted that a resultant picture would not look like a normal picture since the inside of the cave would be just as bright as what surrounds the cave on the outside. However, the intent of the circuit of FIG. 1 is to brighten everything to the fullest dynamic range, so that an observer can see what is happening or contained in a resultant picture.\nFIG. 2B illustrates a possible guide as to what may be done to increase the dynamic range of I(x) of FIG. 2A. As shown, FIG. 2B illustrates a plurality of eight local data regions, each 32 pixels in length, which respectively encompass contiguous 32-pixel long portions of I(x). In each data region the brightest (or largest intensity) and darkest (or smallest intensity) parts of the associated portion of I(x) are determined. The largest and smallest intensity values found in each data region are used to form horizontal ceiling (C) and floor (F) values for each pixel across that region. Thus, the sequence of ceiling values found in the respective data regions forms a ceiling function of X values, designated C(x), while the sequence of floor values found in the respective data regions forms a floor function of X values, designated F(x), for the included portions of I(x). FIG. 2C shows just F(x) the floor function of FIG. 2B. As shown, F(x) is comprised of the segmented sequence of minimum values across the respective regions of FIG. 2B.\nAs described, a first possible operation that can be performed to maximize the dynamic range of I(x) is to subtract F(x) from I(x) to produce the wave form shown in FIG. 2D. It can be seen in FIG. 2D that the amplitude of the signal I(x) of FIG. 2A is substantially reduced, while still substantially retaining the intensity variations of FIG. 2A.\nAs also described, the next possible operation that can be performed is to expand the amplitude of the signal I(x)–F(x) of FIG. 2D by adjusting the gain in the first data region (0–32), second data region (32–64), third data region (64–96) etc.—each gain adjustment being independent of the others-so that the maximum value of the signal in each data region go all the way up to the maximum permissible level or saturation level (Isat) FIG. 2E illustrates a piecewise constant gain function G(x) that would independently adjust the amplitude in each segment or data region up to the saturation level. This piecewise constant gain function would be determined by the value of the saturation intensity (Isat) divided by the difference between the C(x) and F(x) functions.\nFIG. 2F illustrates the result of multiplying the function of FIG. 2D by the gain function of FIG. 2E. This multiplication raises the maximum value in each data region up to the saturation level. However, Vitols et al. describe FIG. 2F as being very unsatisfactory because of the discontinuities or, places where the amplitude of the signal rises or falls vertically. Although the full dynamic range of the signal is obtained in each data region, the resultant signal shown in FIG. 2F is very bumpy. So the overall information that an operator may be looking for in the signal may be totally lost in the discontinuities. If the one-dimensional wave form of FIG. 2F were applied to a two-dimensional picture, a very strong checkerboard pattern would result, which would distort the operator's perception of a picture or image to the point where he probably could not discern what he was looking at. Vitols et al. go on to explain how the discontinuities shown in FIG. 2F can be eliminated by smoothing the signal functions F(x) and G(x) from one data region to another. This smoothing operation is described further in the Vitols et al. patent, but will not be discussed further herein.\nOne of the drawbacks of the Vitols et al. circuit is that it requires a complex sequence of operations to determine the desired adjustment functions for the system. As described, the I-getter (as illustrated in FIG. 12 of Vitols et al.) requires two shift registers for every data cell used in an output picture. In addition, a large number of switching points capable of switching out each of the data cells for separate comparator operations must be used. In addition, during the comparator operations the comparator voltage level to which the incoming pixel signals are compared may be constantly shifting. In addition, as further described in Vitols et al., it takes time for each of the function generators to get sufficient parameter data in before the function generators can start computing the associated correction signals. Furthermore, it is also stated that it takes additional time before each of the function generators can compute its associated correction signals. This combined delay time must be offset by a delay circuit which synchronizes the time of occurrence of the sequence of pixels with the time of generation of the pixel bias correction signals and pixel gain correction signals.\nThe present invention is directed to a circuit that overcomes the foregoing and other problems in the prior art. More specifically, the present invention is directed to an auto black expansion method and apparatus for an image sensor that uses a simplified digital control system and does not require additional shift registers for the pixel signals or continual shifting of input comparator levels during a given field."}
{"text": "Many off-road vehicles, such as earthmoving vehicles (e.g., excavators, skid steer track loaders or multi-terrain track loaders) and agricultural vehicles, include an endless track coupled to a drive system, which may include, for example, one or more drive sprockets. An endless track may help provide traction on surfaces that traditional wheels having a circular cross-section may not be able to adequately grip. For example, a vehicle with an endless track may bridge a greater surface area than a vehicle with two or more circular wheels, which may help the vehicle with the endless track traverse surfaces that the vehicle with circular wheels would not ordinarily be able to traverse, such as uneven or loose surfaces (e.g., gravel or mud).\nIn some cases, an endless track is formed of metal, rubber, such as a molded rubber, or a combination of metal and rubber. The endless track may be molded as a single, integral piece or an assembly of interconnected parts that are coupled together to define an endless track. The endless track typically includes a body with a first, interior side that engages with a drive system and a second, exterior side that includes multiple protrusions (i.e., “traction lugs” or “tread lugs”) extending from one or more common surfaces and defining a tread pattern. The exterior side is configured to engage with the ground, and the tread pattern is typically designed to optimize performance of the endless track. For example, the tread pattern may be designed to optimize performance of an endless track for a high speed skid steer track loader on a wet surface, on a loose surface or otherwise."}
{"text": "1. Field of the Invention\nThis invention is related in general to the wireless computing. In particular, the invention consists of a system for securely booting from and operating from a remote data storage device.\n2. Description of the Prior Art\nMobile computing has become extremely popular with today's computer users. The design of today's mobile computers is a trade-off between performance, size, weight, battery life, and cost. One limiting factor is the cost, weight, size, and power requirements of memory storage devices, such as hard disk drives, for storing operating systems and software applications. However placing a memory storage device within a mobile computer creates a potential for a single-point failure, resulting in the loss of all data in the event of damage or theft and an inability to boot the computer, i.e., load the operating system into working memory.\nA primary function of a local memory storage device is to hold start-up instructions for a computer's boot process. A typical boot process 10 is illustration in the flow chart of FIG. 1. In step 12, the mobile computer is powered-on and a power-on self test (“POST”) is executed. The potential locations for an operating system boot partition are identified in step 14. If an operating system boot partition may be located on a PCI adapter, a basic input/output system (“BIOS”) for the peripheral communications interconnect local bus (“PCI”) adapter is loaded into working memory in step 16. Once the operating system boot partition has been located, the operating system's executable instructions are loaded into working memory (step 16) and the boot process is turned over to the operating system (step 18). Finally, communication standards, such as internet protocols (“IP”), are configured to allow the computer to communicate with remote devices.\nA potential solution to single-point failure is to transfer information to remote memory storage devices such a disk servers or tape libraries. A common method of communicating with these off-system devices is to utilize a communication network. Typically, such communication networks have employed local-area networks (“LANS”), wide-are networks (“WANS”), or storage-area networks (“SANS”). These networks have traditionally utilized networks conforming to IEEE standards, such as IEEE 802.3 (Ethernet). Using the IEEE 802.11b standard, mobile computers may be connected to the base stations through the transmission and reception of radio signals. Additionally, the IEEE 802.11b specification allows for the wireless transmission of approximately 11 million bits per second (“Mbps”) of digital data at indoor distances up to a few hundred feet and outdoor distances up to tens of miles on a 2.4 gigahertz (“GHz”) radio broadcast band.\nThe wireless broadcast communication networks are primarily used for Internet protocol (IP) based communication, but can be used for just about any type of digital communication. The area of effective transmission and reception depends on transmission strength, reception antennae, and line-of-sight obstructions.\nUtilizing these wireless communication networks, mobile computers can easily access software applications from and store data to remote storage devices. However, booting a computer and loading an operating system from a remote device is more problematic. Additionally, providing a secure connection absent a local memory storage device can be difficult.\nIn U.S. Pat. No. 6,141,705, Sanjay Anand et al. disclose offloading processors tasks, including data encryption and message circulation. However, Anand's disclosed invention is not directed to mobile computing and does not discuss booting a device absent a local memory storage device. It would be advantageous to have a system for providing a secure connection between a mobile computer and a remote memory storage device that includes an ability to load boot instructions from a remote memory storage device.\nIn U.S. Pat. No. 5,872,968, Richard I. Knox et al. disclose a data processing network with a client connected to a server. The client issues a boot request (“RPL”) and the server sends bootstrap code to the client. The bootstrap code is loaded into client memory and causes the client to issue a bootstrap request from a second server. The system disclosed by Knox is an attempt to boot a client computer in a mixed network environment. While potentially effective in allowing a boot process to occur in a mixed network protocol environment, Knox's invention utilizes conventional encryption techniques to provide security. It would be advantageous to have a mixed security model environment to resolve a single network unsecured boot process.\nIn U.S. Pat. No. 5,896,499, Mark A. McKelvey discloses an embedded security processor used in conjunction with a main processor to provide security for a computer system that is accessible over a computer network. McKelvey's system includes a security processor, preferably an expansion card plugged into the computer system's data bus, that secures access from one network to another and may optionally include a firewall. However, McKelvey does not teach a method for securing the delivery of a boot image to a mobile computer sans memory storage device. Accordingly, it would be advantageous to utilize a disk management session (“DMS”) for securing a remote boot process using a combination of trusted and untrusted connections."}
{"text": "Field of the Invention\nThe present invention relates to methods for improving game interactions, and more particularly to manipulating game objects within a video.\nDescription of the Related Art\nOnline games allow players to provide input to the game using various input devices. The input devices provide controls and the users use these controls to manipulate game play by maneuvering the various game objects within the game. Online game providers and operators continually seek to provide new ways to provide controls and to improve the input devices in order to make it easy for the users to provide input and to manipulate the game play of the games. As online games are becoming popular, game developers, both newer and existing ones, are introducing new games or evolving existing games to interface with newer or varying control mechanisms for controlling input to the games causing unnecessary complications. It is, therefore, becoming increasingly harder for users to memorize the various controls and/or become familiar with the control mechanisms to manipulate game objects in the games, as these control mechanisms seem to largely vary from one game to the next.\nIt is, therefore, advantageous for the online game operators to seek a more simplified solution for providing control mechanisms to manipulate the game objects during game play so as to keep players interested in the game and enriching the players game experience, because increased playing times and satisfied players often means higher game revenues for the game operators.\nIt is in this context that embodiments arise."}
{"text": "Injection molding systems incorporating so-called “hot runners” for distribution of molten thermoplastic material to one or more injection nozzles are generally known. In a typical hot runner system, the injection nozzles project away from a face of the hot runner system for insertion into a splittable mold formed by a first mold plate referred to as the core plate and a second mold plated referred to as the cavity plate such that the tip of the injection nozzle projects to a mold cavity formed between the mold plate and the cavity plate. Injection of thermoplastic material into this cavity causes a part to be formed conforming generally to the contours of the cavity. Upon separating the mold plate from the cavity plate, the cavity is exposed and the formed, solidified part may be withdrawn. Of course, numerous variations to this basic process may be utilized as desired.\nRegardless of the final arrangement of the hot runner system and associated mold, it may be necessary to periodically service the injection nozzles and supporting equipment within the hot runner system. In previous systems, the hot runner generally has been integrally constructed with the molten material feed system and has not been readily susceptible to rapid replacement. Thus, any maintenance required on the hot runner system has necessitated the removal of the mold plates, followed by the on-site repair of the hot runner system. Such on-site repair removes the injection molding machine from use until the repair is fully completed. Specifically, repairing the hot runner system within the injection molding machine may result in the injection molding machine being out of operation for a period of several days while the repairs are taking place. This period of lost service may correlate to a substantial financial loss as parts cannot be manufactured during that time.\nBy way of example only, and not limitation, one prior system that may be used to provide access to a hot runner system is illustrated and described in U.S. Patent Application Publication US2006/0269649 A1, having a publication date of Nov. 30, 2006. As best understood, in the process disclosed in this reference, the mold plates are closed together and secured so as to form a unitary mold module. The mold plates are thereafter released from the hot runner manifold plate and the mold module is pulled away from the runner and lifted out of place, thereby exposing the hot runner for access during a repair. However, as best understood, there is no mechanism provided for removal of the complex hot runner system. Rather, the hot runner may be required to remain in place during the entire period of repair, thereby rendering the entire injection molding system inoperable during that period."}
{"text": "The present invention relates to a stirring device for liquid material, which circulates and stirs liquid material such as plastic material. More specifically, the invention relates to a stirring device wherein a return pipe (i.e. an outlet pipe) is inserted in a stirring tank and a feeding pipe (i.e. an inlet pipe) is connected to bottom portion of the stirring tank so that liquid material withdrawn from the stirring tank through the feeding pipe is returned to the stirring tank via the return pipe.\nIn a molding device to mold foamed polyurethane products, for example, polyol and isocyanate are put in separate stirring devices and then stirred while circulated between the stirring device and an injection device respectively. Such a stirring device of the prior art, used in a plastic molding device, as shown in FIG. 1 comprises a stirring tank 21, a stirring blade 22 driven to rotate at the center within the stirring tank 21, a return pipe 23 inserted in the stirring tank 21 from above for introducing a circulated liquid material into the tank, and a feeding pipe 24 connected to the bottom portion of the stirring tank 21, so that the upper opening of the stirring tank 21 is sealed, a pressure air pipe (not shown) is connected to the upper side and the stirring is performed under pressure in the stirring tank 21.\nIn this construction, however, since the return pipe 23 with its opening end directed downward is arranged in the stirring tank 21 and the liquid material of large flow rate flows towards the bottom portion of the stirring tank 21a, large convection in the vertical direction is produced within the stirring tank 21 and the liquid surface is swollen significantly, whereby the air content in the liquid material starting material for the foamed polyurethane becomes excessive and cannot be controlled suitably."}
{"text": "A computer application typically provides a printable page to a printing device for printing to a hard copy medium, such as a paper sheet. The printable page is typically provided in the form of a description of the page to be printed, specified using a Page Description Language (PDL), such as Adobe® PDF or Hewlett-Packard® PCL. The PDL provides descriptions of objects to be rendered onto the page in a rendering (or z) order, as opposed to a raster image (i.e. a bitmap of pixel values) of the page to be printed. The page is typically rendered for printing by an object-based graphics system, also known as a Raster Image Processor (RIP). A RIP may also be used to render the page to a display.\nThe printing device receives the description of the page to be rendered and generates an intermediate representation of the page. The printing device then renders the intermediate representation of the page to pixels which are printed to print media, such as paper. In general, an intermediate representation of a page consumes less memory than the raster image representation. Also, in some prior art printing devices, the intermediate representation of the page may be rendered to pixels in real-time, being the rate at which the output device, be it a printer or a display, can reproduce output pixels. Real-time rendering is particularly important for video displays, where animation frame rates must be met to ensure fluidity of motion. Real-time rendering in a printing environment is important to ensure compliance with page throughput rates of the printer.\nThe intermediate page representation is generated by a controlling program which is executed by a controlling processor within the printer device. A pixel rendering apparatus is used to render the intermediate page representation to pixels. The rendered pixels are transferred to a printer engine, such as an electro-photographic engine, which prints the pixel onto the print media.\nNext generation printing systems are expected to operate at a much higher page rate than current printing systems. This is in addition to an increase in device resolution, graphics complexity, and the number of print features supported.\nIn computing, there is a trend towards achieving higher performance through the use of multi-processor and multi-core architectures. These architectures allow a number of threads to execute in parallel, each on a corresponding processor, thereby reducing the overall time it takes to complete a task. However, in order to take advantage of such parallelism, the task must be broken down into largely independent sub-tasks that can be executed in parallel. This is difficult to achieve for many tasks, including many of the operations performed by a RIP.\nThe RIP process consumes a large proportion of time within a printing system. It is therefore desirable that the RIP process be accelerated through the use of multi-threading on a multi-core or multi-processor device. One method of performing the RIP process in parallel is by sub-dividing the output pixel space into regions, and performing the RIP process for each region in its own thread.\nOne method of doing this is to assign each page object to be rasterized to each of the regions that the page object overlaps or is present in. The regions can then be rasterized in parallel, with a separate thread processing each region. Each thread can only process those objects that overlap the region that is being rendered. However, because many objects may overlap many regions, many objects must be scan-converted multiple times, once for each region that the object overlaps. It is desirable that this duplication of processing be removed.\nAnother method of sub-dividing the output pixel space into regions, which removes this duplication of object scan-conversion, is to split each page object prior to scan-conversion, according to the object's intersections with region boundaries. A number of sub-objects are produced, each of which is present in (or overlaps) only a single region. Each sub-object therefore only needs to be scan-converted once. However, the splitting of high-level graphic objects (which normally consist of primitives such as lines and Bezier splines) is complex, and can result in mathematical errors which affect the quality of the rasterized output.\nIt is desirable that one or more these problems with the prior art be resolved, or at least ameliorated, while still allowing the intermediate representation of each region of the page to be generated in parallel."}
{"text": "This invention relates to a carbonaceous granular heat insulator for use in a high-temperature oven in a non-oxidizing atmosphere and a process for preparing the same.\nA carbon material is excellent in a high-temperature stability and a chemical stability in a non-oxidizing atmosphere and, therefore, is used as a refractory material and a heat insulator in a high-temperature oven.\nIn general, an insulating performance required of an insulating material, i.e., thermal conductivity, greatly varies depending on the porosity of the insulating material. As the porosity increases, the thermal conductivity decreases while the insulating performance is improved.\nExamples of the carbonaceous heat insulator which has been put to practical use include a porous carbonaceous material prepared by carbonizing a foamed resin through baking or a bulky carbon molding prepared by impregnating a felt made of a carbon fiber with a thermosetting resin, molding the impregnated felt, curing the resulting molding and carbonizing the cured molding through baking. All of these carbonaceous heat insulators are foamed carbon material and, therefore, take advantage of an insulating function of large voids present therein.\nHowever, the raw materials of these foamed carbon materials are disadvantageous in that the applications thereof are limited because the cost in high due to expensive raw materials used and complicated steps of production.\nCarbon black comprises an aggregate having a chain structure formed by fusion and secondary aggregation of very fine spherical fundamental particles having a diameter of 10 to 200 nm, and many voids are present among the aggregate.\nAlthough the surface layer and its vicinity of the fundamental particle comprise a carbon hexagonal network, the fundamental particle comprises amorphous carbon on the whole. The carbon black itself, therefore, has an excellent insulating property and is advantageously inexpensive, which makes it possible to put the carbon black into practical use, e.g., as an insulating filler of a high-temperature oven.\nHowever, since the carbon black is a finely divided powder comprising an aggregate of fine fundamental particles, it tends to scatter during charging of an oven therewith and discharging it from the oven, which brings about environmental pollution. Further, it is poor in fluidity, which makes it very difficult to handle. In general, carbon black to be compounded with rubber is granulated for the purpose of controlling the environmental pollution and improving the fluidity. However, pellets prepared by granulation are fragile and do not have a pellet hardness sufficient to prevent the occurrence of attrition loss and breakage of the granules. Further, the pellets tend to break into finely divided powder, which raises problems such as scattering of the powder and lowering in fluidity."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the packaging of electronic components. More particularly, the present invention relates to an electronic component package.\n2. Description of the Related Art\nImage sensors are well known to those of skill in the art. An image sensor included an active area, which was responsive to electromagnetic radiation. The image sensor was incorporated into an image sensor package, which protected the image sensor from dust and moisture.\nIn one image sensor package, the image sensor was located within an enclosure, which included a window. Electromagnetic radiation passed through the window and struck the active area of the image sensor, which responded to the electromagnetic radiation. The enclosure also served to protect the image sensor from dust and moisture. Although the enclosure around the image sensor was reliable and effective, fabrication of the enclosure was relatively complex, time-consuming, and thus expensive.\nIn accordance with one embodiment of the present invention, a method of forming an image sensor package includes coupling an image sensor, i.e., an electronic component, to a substrate. Bond pads of the image sensor are electrically coupled to interior traces on the substrate with bond wires. A first optically curable material is applied to enclose the bond wires. A second optically curable material is applied between a lid and the substrate. The first and second optically curable materials are cured with a first electromagnetic radiation such as ultraviolet radiation. The substrate, lid, and cured second optically curable material form a low cost enclosure around the image sensor.\nIn one embodiment, the first and second optically curable materials are optically entirely curable materials. In accordance with this embodiment, the first and second optically curable materials are cured immediately as a result of being irradiated with the first electromagnetic radiation. More particularly, the first and second optically curable materials are entirely cured while being irradiated with the first electromagnetic radiation.\nIn an alternative embodiment, the first and second optically curable materials are optically activated materials. In accordance with this embodiment, the first and second optically curable materials cure over a period of time as a result of being irradiated with the first electromagnetic radiation. More particularly, the first and second optically curable materials cure initiate while being irradiated with the first electromagnetic radiation and then entirely cure over a period of time without further irradiation.\nIn one embodiment, after being cure initiated, the first and second optically curable materials entirely cure at room temperature, i.e., without heating. Alternatively, after being cure initiated, the first and second optically curable materials entirely cure at elevated temperature, e.g., are heated to reduce the curing time.\nAdvantageously, the first and second optically curable materials are cured rapidly, e.g., in a matter of seconds. This is in stark contrast to conventional thermally curable materials, which often had a cure time of several hours. By curing the first and second optically curable materials rapidly, the time required to fabricate the image sensor package is minimized. This, in turn, minimizes the fabrication cost of the image sensor package.\nIn one embodiment, the first and second optically curable materials are cured without heating the image sensor package to any appreciable extent. This is in stark contrast to conventional thermally curable materials, which required heating for several hours.\nAccordingly, use of the first and second optically curable materials is well suited for applications in which it is important to keep the image sensor and any other temperature sensitive components of the image sensor package at low temperature, for example, to prevent degradation or damage to any color filters and/or other temperature sensitive materials of the image sensor.\nFurther, since the first and second optically curable materials are cured without heating, the requirement for cure ovens, which were necessary to cure conventional thermally curable materials, is eliminated further minimizing the fabrication cost of the image sensor package.\nAlso in accordance with one embodiment of the present invention, an image sensor package includes a substrate and an image sensor coupled to the substrate. An interior trace is coupled to the substrate. A bond wire electrically couples a bond pad of the image sensor to the interior trace. An encapsulant encloses the bond wire, the encapsulant being formed of a first optically curable material that has been cured.\nThe image sensor package further includes a lid. A lid adhesive couples the lid to the substrate, the lid adhesive being formed of a second optically curable material that has been cured."}
{"text": "1. Field of the Invention\nThe present invention relates to generally to phosphors and more specifically to the production of nanoparticulate phosphors.\n2. Description of the Background Art\nThere is a resurgence in phosphor research due to the recent interest in developing emissive flat panel displays, e.g. full-color low-voltage field emitter displays (FEDs), large area plasma displays (PDs) and electroluminescent displays (ELDs). For example; to make low-voltage FEDs. decreasing the required electron energy by two orders of magnitude from the 30 KV range of CRTs while maintaining significant light output is quite challenging. In low-voltage FEDS, the incident electron mean free path in the phosphor is extremely short, and some researchers feel that the photon generation process should occur reasonably close to the electron-hole pair creation point to avoid non-radiative recombination at nearby defects. Furthermore, there is the belief that present day micron size phosphors are surrounded by a dead layer of defective material whose thickness is comparable to or greater than the mean free path of low energy electrons used in low voltage FEDs. The defects in this surface layer, which are produced, for example, by ball milling and other harsh surface (and bulk) defect producing processes, are believed to be non-radiative recombination centers. These result in low optical output when these phosphors are bombarded by low energy electrons.\nResearch interest in the general field of nanoscale materials has been growing fast. New methods of developing nanoparticles are being developed, and their novel physical properties are being studied. For example, in Yang et al. U.S. Ser. No. 08/841,957 filed Apr. 7, 1997. Now U.S. Pat. No. 6,106,609 nanocrystalline semiconductor particles are produced within the mesoporous cavities of a bicontinuous cubic phase. The cavities contain at least one of the reactants required to form the semiconductor compound. Any other required reactants are diffused into the mesoporous cavities.\nAccordingly, it is an object of the present invention to provide phosphors having a high light output level.\nIt is another object of the present invention to provide nanocrystalline phosphors that have no or few bulk defects.\nThese and additional objects of the invention are accomplished by forming doped nanocrystalline semiconductor particles within the mesoporous cavities of a bicontinuous cubic phase. The cavities contain at least one of the reactants required to form the semiconductor compound. Any other required reactants are diffused into the mesoporous cavities. Doped nanocrystals are produced by diffusing dopant or activation ions into the mesoporous cavities along with the other reactants or by including dopant or activation ions within the mesoporous cavities of the BCP under conditions such that the nanoparticles are doped with an average of one or less dopant ions per particle. For the purposes of the present specification and claims, the words xe2x80x9cdopantxe2x80x9d and xe2x80x9cactivatorxe2x80x9d are used interchangeably to describe the ion that when placed within the interior of a nanocrystalline particle, provides the particles with phosphorescence."}
{"text": "For those who spend periods of time sitting in a bed, seating, furniture, sofa or even on a floor and/or against a wall, it is important to provide proper support to particular regions of the body to avoid discomfort, fatigue and/or body deformation and musculo-skeleton strains. Ordinary cushions and pillows fail to provide necessary support and structural integrity, which requires a chair-like bed lounge typically with at least a back portion with two attached arm-rests, one at each side, and preferably a neck support/head-rest. To accommodate the needs of different individuals with regard to size, weight, age, activity, etc., the Bed Lounge is preferably provided with several user adjustment capabilities, e.g. relating to tilt-back, arm-rest and seat-back spacing, lower back support, neck-rest height and inclination."}
{"text": "The invention relates to the field of interleaving and deinterleaving of data.\nIn the field of communications, forward error correction (FEC) coding is typically applied to blocks of data that are to be transmitted and this type of coding normally allows successful recovery from errors that affect relatively short parts of a received data block. Interleaving is a technique that is commonly used to reduce the chance that an error will affect a relatively long part of a received data block, as will now be explained.\nPrior to transmission, the data items within an FEC encoded data block can be shuffled into a different order. This shuffling is referred to as interleaving and swaps the data items from a first pre-defined order to a second pre-defined order. After reception, the data items of the block are shuffled back into their original order. This shuffling is referred to as deinterleaving and swaps the data items from the second pre-defined order to the first pre-defined order. If interference during transmission causes an error over a part of the received version of the interleaved block, then the deinterleaving process distributes the error to various locations within the block. That is to say, after deinterleaving, the error is less likely to affect a contiguous part of the received block that is of sufficient length to impede recovery from the (dispersed) error by applying FEC decoding.\nRecent wireless communications standards such as the Long Term Evolution (LTE) project by the Third Generation Partnership Project (3GPP) use large transport blocks. A channel interleaver/deinterleaver for such blocks will require considerable amounts of memory storage. A “single buffer” approach of writing a data block into an addressable memory in an interleaved order and then subsequently reading the data block out of the memory in a deinterleaved order may be undesirably slow, particularly where the block size is large. A “double buffer” approach of writing a data block into one addressable memory in an interleaved order whilst reading another data block out of another addressable memory in a deinterleaved order is faster but is costly in terms of silicon area."}
{"text": "1. Field of the Invention\nThe instant invention relates in general to ramps used for ascending and descending an object or platform and more particularly to a ramp which is lightweight and foldable for ease in transportation and storage.\n2. Description of the Related Art\nThe loading and unloading of utility off-the-road vehicles, lawn and garden equipment, and other apparatus into trucks or vans has presented a problem for some time. Many alternatives have been proposed for addressing this need. Unfortunately, the proposed solutions usually involve ramps which are permanently or semipermanently attached to the truck or van and hence not portable. These prior art ramps are unnecessarily heavy resulting in reduced gas mileage of the truck or van and difficulty in positioning the ramp for use. Further, many of the prior art ramps are bulky resulting in great difficulty in handling and storing the ramp.\nOne example of such a prior art ramp is disclosed in U.S. Pat. No. 3,976,209. That ramp utilizes a plurality of pivotable legs to directly support a folding hinge. The legs provide structural support for the ramp without the requirement of a vertical ramp to ground support. However, the legs are rigid and cannot be adjusted to carry varying weight loads. A further example of a prior art ramp is discussed in U.S. Pat. No. 1,115,775. This ramp includes a vertical support member at the hinge. However, the support member is rigid which precludes folding the ramp into a compact state. Additionally, the support rather, it simply connects the two plank members forming the ramp.\nKnown prior art ramps which are foldable, only hinge in one direction, and therefore the compactness of the ramps is limited. Furthermore, no prior art ramps are known which are adjustable to limited in the size and type of objects which can be ascended with the ramp.\nMoreover, none of the known prior art ramps are light enough to be easily transported, hence alternate uses for the ramp are limited.\nIt is an object of the present invention to provide a portable, lightweight, and compact ramp for loading or unloading utility vehicles or other objects into or out of a truck, car, van or other vehicle.\nIt is a further object of the invention to provide a lightweight ramp which is foldable in more than one direction to make the ramp more compact for storage and transportation and to provide adjustability of the lateral dimension of the ramp so as to accomodate varying sizes of objects or platforms which are to be ascended.\nIt is still a further object of the present invention to provide a lightweight, compact ramp having a load support system which can be selectively adjusted to accomodate varying loads placed on the ramp.\nIt is still a further object of the present invention to provide a lightweight, compact ramp which is easily transported and handled so that the ramp may be used for a multiplicity of purposes.\nIt is still a further object of the invention to provide a portable, lightweight ramp having an internal load supporting mechanism which does not require contact with the ground to provide support to the ramp.\nAdditional objects and advantages of the invention will be set forth in the description which follows and in part will be obvious from the description or may be learned by practice of the invention. The objects and advantages of the invention may be realized and attained by means of the instrumentalities and combinations particularly pointed out in the appended claims."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a semiconductor memory device and a method of manufacturing the same. More particularly, the invention relates to a semiconductor memory device having an etch stopper formed in a cell region and a method of manufacturing the semiconductor memory device.\nA claim of priority is made to Korean Patent Application No. 10-2005-0015371, filed on Feb. 24, 2005, the disclosure of which is hereby incorporated by reference in its entirety.\n2. Description of Related Art\nAs the integration density of semiconductor memory devices increases, contact holes formed in the semiconductor devices become increasingly narrow. As a result, the aspect ratio of the contact holes tends to increase, and therefore contact holes and associated contacts are often required to be formed using multi-step processes.\nThe integration density of dynamic random access memory (DRAM) devices is typically very high. Because of this, active regions between adjacent gate electrodes in DRAMs are generally so narrow that it is difficult to form contacts connected to the active regions. One approach to forming contacts in narrow regions is to form a hard mask layer and spacers over and around the gate electrodes before forming self aligned contact pads between the gate electrodes. Then, bitline contacts or storage node contacts are formed over the self-aligned contact pads. This approach is often referred to as “self aligned contact technology.”\nTypically, the self-aligned contacts have a larger diameter than the bitline contacts or the storage node contacts. However, as the integration density of DRAM devices continues to increase, the diameter of the self-aligned contact pads reduces and approaches the diameter of the bitline and storage node contacts. As a result, slight misalignments between the bitline or storage node contacts and node can cause serious problems in the manufacture of DRAM devices. Illustrations of some common problems are described below in relation to FIG. 1.\nReferring to FIG. 1, a conventional semiconductor memory device comprises a semiconductor substrate 10 in which a cell region “A” and a core and peripheral region “B” are defined. A plurality of isolation layers 15 defining active regions 10a are formed in semiconductor substrate 10, and junction regions 25a and 25b are formed between isolation layers 15 in respective cell and core and peripheral regions “A” and “B”.\nA first insulating interlayer 30 is formed over semiconductor substrate 10 and a second insulating interlayer 40 is formed over first insulating interlayer 30.\nIn cell region “A”, a plurality of self-aligned contact pads 35 are formed in first insulating interlayer 30 over junction regions 25a, and a plurality of bitline contact holes 45 are formed in second insulating interlayer 40 to expose self-aligned contact pads 35.\nIn core and peripheral region “B”, a gate electrode structure 20 is formed in first insulating layer 30 on side junction region 25b and an interconnection contact hole 50 is formed in first and second insulating interlayer 30 and 40 to expose side junction region 25b. Each gate electrode structure 20 comprises a gate insulating layer 16 formed on semiconductor substrate 10, a gate electrode material 17 formed on gate insulating layer 16, a hard mask layer 18 formed over gate electrode material 17, and spacers 19 formed on the sidewall of gate insulating layer 16, gate electrode material 17, and hard mask layer 18.\nTypically, the respective diameters of self-aligned contact pads 35 and bitline contact holes 45 are slightly different from each other, or they are slightly misaligned with each other. As a result, portions of first insulating interlayer 30 adjacent to self-aligned contact pads 35 may be removed when bitline contact holes 45 are formed. Due to the partial loss of first insulating interlayer 30, active regions 10a may also be partially removed, causing leakage current in the semiconductor memory device.\nStorage node contact holes (not shown) are typically formed after bitlines are formed in bitline contact holes 45. The storage node contact holes are generally formed between the bitlines to contact self-aligned contact pads 35. However, the storage node contact holes are almost inevitably misaligned with the self-aligned contact pads. As a result, the first insulating interlayer 30 is almost inevitably etched when the storage node contact holes are formed, causing further damage to active regions 10a. \nInterconnection contact hole 50 is typically formed in core and peripheral region “B” at the same time when bit line contact holes 45 are formed in cell region “A”. Unfortunately, however, as the integration density of core and peripheral region “B” increases, interconnection contact hole 50 tends to expose more than just junction region 25b. For instance, interconnection contact hole 50 may expose parts of gate electrode structure 20.\nConsequently, spacers 19 and hard mask layer 18 may be partially lost during the formation of interconnection contact hole 50. As a result, gate electrode material 17 may be exposed and may be short-circuited by a conductive material formed later in interconnection contact hole 50.\nTo solve these problems, some conventional memory devices interpose etch stoppers between first insulating layer 30 and second insulating interlayer 40 to prevent first insulating interlayer 30 from being etched when bitline contact holes 45 are formed. A conventional memory device employing this technique is disclosed, for example, in U.S. Pat. No. 6,787,906 (the '906 patent).\nIn the '906 patent, etch stoppers are formed on first insulating layer 30. Accordingly, bitline contact holes 45 and storage node contact holes can be formed without further etching first insulating layer 30. However, the etch stoppers are also formed in core and peripheral region “B”, making it difficult to form interconnection contact hole 50. Since interconnection contact hole 50 penetrates first and second insulating layers 30 and 40, if the etch stoppers are interposed between first and second insulating interlayers 30 and 40, it is difficult to expose junction region 25b. "}
{"text": "1. Technical Field of the Invention\nThe invention relates generally to communication systems; and, more particularly, it relates to beamforming of wireless communications within such communication systems.\n2. Description of Related Art\nCommunication systems are known to support wireless and wire lined communications between wireless and/or wire lined communication devices. Such communication systems range from national and/or international cellular telephone systems to the Internet to point-to-point in-home wireless networks. Each type of communication system is constructed, and hence operates, in accordance with one or more communication standards. For instance, wireless communication systems may operate in accordance with one or more standards including, but not limited to, IEEE 802.11, Bluetooth, advanced mobile phone services (AMPS), digital AMPS, global system for mobile communications (GSM), code division multiple access (CDMA), local multi-point distribution systems (LMDS), multi-channel-multi-point distribution systems (MMDS), radio frequency identification (RFID), Enhanced Data rates for GSM Evolution (EDGE), General Packet Radio Service (GPRS), and/or variations thereof.\nDepending on the type of wireless communication system, a wireless communication device, such as a cellular telephone, two-way radio, personal digital assistant (PDA), personal computer (PC), laptop computer, home entertainment equipment, RFID reader, RFID tag, et cetera communicates directly or indirectly with other wireless communication devices. For direct communications (also known as point-to-point communications), the participating wireless communication devices tune their receivers and transmitters to the same channel or channels (e.g., one of the plurality of radio frequency (RF) carriers of the wireless communication system or a particular RF frequency for some systems) and communicate over that channel(s). For indirect wireless communications, each wireless communication device communicates directly with an associated base station (e.g., for cellular services) and/or an associated access point (e.g., for an in-home or in-building wireless network) via an assigned channel. To complete a communication connection between the wireless communication devices, the associated base stations and/or associated access points communicate with each other directly, via a system controller, via the public switch telephone network, via the Internet, and/or via some other wide area network.\nFor each wireless communication device to participate in wireless communications, it includes a built-in radio transceiver (i.e., receiver and transmitter) or is coupled to an associated radio transceiver (e.g., a station for in-home and/or in-building wireless communication networks, RF modem, etc.). As is known, the receiver is coupled to an antenna and includes a low noise amplifier, one or more intermediate frequency stages, a filtering stage, and a data recovery stage. The low noise amplifier receives inbound RF signals via the antenna and amplifies then. The one or more intermediate frequency stages mix the amplified RF signals with one or more local oscillations to convert the amplified RF signal into baseband signals or intermediate frequency (IF) signals. The filtering stage filters the baseband signals or the IF signals to attenuate unwanted out of band signals to produce filtered signals. The data recovery stage recovers raw data from the filtered signals in accordance with the particular wireless communication standard.\nAs is also known, the transmitter includes a data modulation stage, one or more intermediate frequency stages, and a power amplifier. The data modulation stage converts raw data into baseband signals in accordance with a particular wireless communication standard. The one or more intermediate frequency stages mix the baseband signals with one or more local oscillations to produce RF signals. The power amplifier amplifies the RF signals prior to transmission via an antenna.\nWhile transmitters generally include a data modulation stage, one or more IF stages, and a power amplifier, the particular implementation of these elements is dependent upon the data modulation scheme of the standard being supported by the transceiver. For example, if the baseband modulation scheme is Gaussian Minimum Shift Keying (GMSK), the data modulation stage functions to convert digital words into quadrature modulation symbols, which have a constant amplitude and varying phases. The IF stage includes a phase locked loop (PLL) that generates an oscillation at a desired RF frequency, which is modulated based on the varying phases produced by the data modulation stage. The phase modulated RF signal is then amplified by the power amplifier in accordance with a transmit power level setting to produce a phase modulated RF signal.\nAs another example, if the data modulation scheme is 8-PSK (phase shift keying), the data modulation stage functions to convert digital words into symbols having varying amplitudes and varying phases. The IF stage includes a phase locked loop (PLL) that generates an oscillation at a desired RF frequency, which is modulated based on the varying phases produced by the data modulation stage. The phase modulated RF signal is then amplified by the power amplifier in accordance with the varying amplitudes to produce a phase and amplitude modulated RF signal.\nAs yet another example, if the data modulation scheme is x-QAM (16, 64, 128, 256 quadrature amplitude modulation), the data modulation stage functions to convert digital words into Cartesian coordinate symbols (e.g., having an in-phase signal component and a quadrature signal component). The IF stage includes mixers that mix the in-phase signal component with an in-phase local oscillation and mix the quadrature signal component with a quadrature local oscillation to produce two mixed signals. The mixed signals are summed together and filtered to produce an RF signal that is subsequently amplified by a power amplifier."}
{"text": "1. Field of the Invention\nThe present invention relates to a plasma display panel and the manufacturing method thereof, more particularly to the partition wall structure of the panel and the manufacturing method thereof.\n2. Description of the Prior Art\nThe rib of the plasma display panel (referred to PDP in the following) commonly has a stripe-shaped structure. However, the grid-mesh rib structure is also used at present, for example, the one disclosed in the U.S. Pat. No. 5,701,056 by NEC. The structure disclosed by NEC forms stripe-shaped ribs on the back substrate of the PDP and forms grid-mesh-shaped ribs on the front substrate of the PDP, then assembles the front and back substrates, as shown in FIG. 1. The structure disclosed by NEC has the following four disadvantages:\nSince the front substrate has an additional rib manufacturing process in the NEC structure, the cost is relatively high.\nWhen assembling the front and the back substrates, the high aligning precision of the two substrates is strictly required; this deepens the difficulty of the manufacturing process.\nTo ensure that the front and the back substrate are precisely aligned, increasing the width of the rib of the front and the back substrates is often required. Hence the opening rate of the PDP is compromised.\nDue to the width of the rib, the effective area of the coating fluorescent body becomes smaller."}
{"text": "Scaling of the gate dielectric is a challenge in improving performance of advanced field effect transistors. In a field effect transistor employing a silicon oxide based gate dielectric, the leakage current through the gate dielectric increases exponentially with the decrease in the thickness of the gate dielectric. Such devices typically become too leaky to provide high performance at or below the thickness of 1.1 nm for the silicon oxide gate dielectric.\nHigh-k gate dielectric provides a way of scaling down the thickness of the gate dielectric without an excessive increase in the gate leakage current. However, high-k gate dielectric materials are prone to a change in the effective oxide thickness (EOT) because high-k gate dielectric materials react with oxygen that diffuses through the gate electrode or gate spacers. Regrowth of a silicon oxide interfacial layer between a silicon substrate and the high-k gate dielectric during high-temperature processing steps is a major obstacle to successful effective oxide thickness scaling. Particularly, typical stacks of a high-k gate dielectric and a metal gate is known to be susceptible to a high temperature anneal in an oxygen ambient. Such a high temperature anneal in oxygen ambient results in regrowth of the silicon oxide interfacial layer and produces instability of the threshold voltage of field effect transistors."}
{"text": "Our co-pending application number WO-A-02065973 discloses a method of treatment of lymphodema and leg ulcers and a prophylactic treatment for deep vein thrombosis (DVT) employing mechanical vibrations, particularly cycloidal vibration, employing a vibration pad. Cycloidal vibration is a small amplitude, 0.1 and 0.5 mm, low frequency, 15 to 75 HZ, vibration that produces motion in three different directions, each of these directions will be at different points in its cycle. It is the out of “phase” relationship which gives rise to the term cycloid vibration. GB-A-2096899 and U.S. Pat. No. 3,019,785 disclose a vibration pad device comprising a motor mounted in a frame, the frame extending into a pad and the motor driving an eccentric weight that causes cycloid vibration of the pad. Cycloidal vibration can be administered by means of integration of the mechanism into static products such as a portable pad as disclosed in GB-A-2096899, but equally it can be incorporated in a mattress, of a therapy couch, for example.\nU.S. Pat. No. 2,006,247601 relates to the treatment of cellulitis. Cellulitis is a common skin infection. In 2002 to 2003, in the UK, there were nearly 60,000 recorded admissions into hospital. Each admission can take on average 10 days to treat (2), accounting for up to six hundred thousand-bed days per annum. Most commonly affecting the lower limbs, cellulitis is an acute infection of the skin and subcutaneous tissues, characterised by: local heat, redness, pain, erythematous tissue and swelling. It is commonly caused by the bacteria streptococci and is associated with, or can be a consequence of, lower limb swelling/oedema. This can be due to a mix of any of the following: leg oedema, venous hypertension, lymphodema, chronic ulceration and immobility.\nU.S. Pat. No. 2,006,247601 provides a method of treatment of cellulitis comprising the steps of administering one or more antibiotics and applying a vibration pad to the region of the skin affected by cellulitis and submitting the pad to cycloid vibration for a period of at least 30 minutes at least once per day until the infection diminishes. The vibrations are believed to assist transport of the antibiotics to the site of infection so that they have their effect more rapidly and completely.\nLeg ulcers, lymphodema and cellulitis all result, to a greater or lesser extent, in exudation from the skin of potentially infectious material. WO-A-02065973 proposes the use of a cover for the vibration pad, so that exudates can be isolated from the pad and any weeping or bleeding of wounds can be absorbed by the cover and not infect the surface of the pad. However, there is a strong tendency in all environments to reuse apparatus and in today's environmentally-conscious society the temptation not to waste disposable medical supplies can sometimes lead to risks being taken. In fact, the danger of cross-infection is far more costly, even just in environmental terms, than a strict adherence to a single-use policy regarding medical products, and this includes covers of vibration pads.\nIt is an object of the present invention to provide a cover including means to substantially guarantee single use thereof, or at least to make it difficult to reuse."}
{"text": "1. Field of the Invention\nThis invention relates to a culture apparatus, more particularly to a plant culture apparatus.\n2. Description of the Related Art\nWith promotion of plant cultivation techniques, a number of ornamental plants can be currently planted in a closed plant culture apparatus. Generally, as shown in FIG. 1, the conventional closed plant culture apparatus 1 includes a transparent culture container 11 and a cover 12 to air-tightly close the culture container 11. A culture medium 10 such as agar is disposed in the culture container 11, and roots 21 of the ornamental plant 2 are planted in the culture medium 10. Hence, the plant 2 is able to obtain nourishment from the culture medium 10 and carry out photosynthesis utilizing the light penetrating through the transparent culture container 11.\nHowever, the inner surface of the culture container 11 is usually adhered with moisture produced from volatilization of the culture medium 10 and photosynthesis of the plant 2. The inner surface of the culture container 11 is then fogged by the moisture adhered thereto, so that the plant 2 cannot be seen clearly from the outside, and that the ornamental value of the plant culture apparatus 1 is decreased. Additionally, the nourishment included in the culture medium 10 is limited, and the plant 2 tends to wilt due to lack of nourishment after being planted in the plant culture apparatus 1 for a period of time. Hence, the period for the plant 2 to be admired is relatively short."}
{"text": "Field of the Invention\nThe present invention relates to the field of consumer accounts, business accounts, and other account types. More particularly, the present invention relates to a method and system for linking source accounts of, for example, demand deposit accounts (DDAs), credit cards, debit cards, stored value cards, ATM cards, loyalty cards, membership cards, and identification cards to a single account platform. This platform and all its linked source accounts can be accessed by a number of potential remote access devices.\nDescription of the Related Art\nThe ever-bulging wallet is a guaranteed phenomenon in modern-day society. As one traverses through life, one will undoubtedly accumulate numerous things for the wallet. One of those things is the ubiquitous “plastic” in the form of, for example, a credit card, debit card, stored value card, ATM card, phone card, or all of the above. The average consumer has multiple bankcards, private label cards (such as store credit cards, oil/gas credit cards), stored value cards, loyalty cards, membership cards, and identification cards that he/she may carry around. The proliferation of these cards and their associated accounts is adding complexity and inconvenience to consumers' financial lives."}
{"text": "Chromatographic methods generally are used to separate and/or purify molecules of interest such as proteins, nucleic acids, and polysaccharides from a mixture. Affinity chromatography specifically involves passing the mixture over a matrix having a ligand specific (i.e. a specific binding partner) for the molecule of interest bound to it. Upon contacting the ligand, the molecule of interest is bound to the matrix and is therefore retained from the mixture. Affinity chromatography provides certain advantages over other types of chromatography. For example, affinity chromatography provides a purification method that can isolate a target protein from a mixture of the target protein and other biomolecules in a single step in high yield.\nDespite the advantages of current affinity chromatography devices, there exists a need in the art for a chromatography device that can be used at shorter residence times than conventional devices while providing the same binding capacity or better binding capacities than current offerings and that is re-useable."}
{"text": "This invention relates to a reflow soldering apparatus and method, particularly, comprising a pressure reducing chamber capable of reducing gas at the soldering part when soldering."}
{"text": "The present invention relates to the field of slot-dipole antenna elements, and in particular to the use of such antenna elements in arrays for aerospace applications.\nAntennas are required for many aerospace applications, such as electronically scanned arrays for aircraft or satellite radar and communications systems or missile tracking, telemetry, and seeker antennas. The radiating elements used in such applications must conform to the surface of the vehicle carrying the antennas and must be both lightweight and capable of being manufactured relatively inexpensively and accurately using printed circuit technology.\nModern surveillance radars also require a wide signal bandwidth for scanning. The pattern beamwidth appropriate for wide angle scanning may also require dual orthogonal senses of polarization. Some commonly-used printed circuit elements for conformal array applications include a microstrip patch, a printed circuit dipole, and stripline-fed, cavity-backed slots. These elements usually have a narrow bandwidth, typically around three percent (3%), which limits their utility. Other commonly used radiating apertures for antenna arrays consist of metallic rectangular or circular waveguides or cavities. These waveguides or cavities, however, are expensive to manufacture and are prohibitively heavy for airborne applications."}
{"text": "1. Field of the Invention\nThe present invention relates to a power storage device that can be reduced in weight and a method for manufacturing the power storage device.\n2. Description of the Related Art\nPortable information terminals such as mobile phones and portable game machines are equipped with power storage devices. As such power storage devices, lithium secondary batteries have been used.\nAlthough portable information terminals have been reduced in weight, the reduction has been unsatisfactory. Portable information terminals equipped with light-weight power storage devices can reduce the weight of the portable information terminals.\nA solid-state secondary battery is a kind of power storage device (see Patent Document 1). The solid-state secondary battery includes a solid electrolyte such as an inorganic solid electrolyte and an organic solid electrolyte.\nIn the technical field of display devices which is different from that of power storage devices, a method for transferring a layer to be separated to a plastic substrate is described (Patent Document 2). Patent Document 2 discloses the following method: a nitride layer is formed over a substrate, an oxide layer is formed over the nitride layer, and film stress of the oxide layer is made different from film stress of the nitride layer to enable peeling by a physical means."}
{"text": "Armorers and gunsmiths are people whose duties include repairing and configuring firearms. One class of tools that have been commonly used is armorer's blocks and magazine posts. Armorer's blocks, also called bench blocks, are typically items that can be laid on a bench and have cavities, openings, slots, and holes that ease various tasks such as driving pins into and out of a firearm or steadying a workpiece. Armorers and gunsmiths typically have and use a variety of armorers block in their day-today labors. Armorers and gunsmiths working on semiautomatic firearms also use magazine posts.\nMagazine posts are posts that can hold a semiautomatic firearm. For example, a bench vise can hold the bottom portion of a magazine post while the top portion of the magazine post is in the magazine well of a semiautomatic firearm. One popular magazine post is configured to slide into the magazine well of an AR-15/M-4 style rifle/carbine and to engage the magazine release mechanism. The lower portion of the AR-15/M-4 magazine post can be clamped in a vise to solidly hold the firearm's lower receiver in a fixed position.\nMagazine posts and armorer's blocks are prized tools in a gunsmith's toolbox because they significantly improve a person's ability to work on a firearm. The currently available tools, however, are not ideal for many modern day semiautomatic handguns or for modern day workbenches. Systems and methods for fixedly holding modern day semiautomatic handguns are needed."}
{"text": "1. Field of Invention\nThis invention relates to accessing data networks by using voice communication devices.\n2. Description of Related Art\nData networks such as the Internet have gained widespread acceptance as a desirable communication medium as well as a source of information. Conventionally, data networks are accessed via devices such as personal computers. Such devices provide visual outputs through a video display and data entry through a keyboard, for example. Additionally, personal computers provide a local storage medium for retrieving data from the data network and storing user preferences, for example. However, potential users that do not have access to a personal computer, for example, are prevented from taking advantage of the benefits provided through a data network. Thus, new technology is needed to make available the data network in such circumstances."}
{"text": "The standard in the United States for the transmission of digital television signals is known as 8 VSB data (ATSC Digital Television Standard A/53). This 8 VSB data has a constellation consisting of eight possible symbol levels. In a VSB system, the eight possible symbol levels are all in the same phase. In a QAM system, however, the symbols are transmitted in phase quadrature relationship.\nThe standard referred to above specifies the formatting and modulation of digital video and audio data. The transmitted data is in the form of symbols with each symbol representing two bits of data that are trellis encoded into three bits of trellis encoded data. Each three bits of trellis encoded data are mapped into a symbol having a corresponding one of eight levels. Reed/Solomon encoding and interleaving are also provided to increase the robustness of the transmitted information.\nAuxiliary data (data other than digital video or audio data) are also permitted to be transmitted in a digital television channel. These data are formatted and modulated according to the standard in the same manner as video and audio data. Receivers made in accordance with the 8 VSB standard are able to read packet identifications (PIDs) which allow the receivers to differentiate between audio, video, and auxiliary data.\nHowever, while the robustness of the transmitted digital television signals is sufficient for digital television reception, this robustness may not be sufficient for the transmission of auxiliary data, particularly where the auxiliary data are critical. Accordingly, one of the applications of the present invention is the transmission of auxiliary data in a VSB format with outer encoding for added robustness. The auxiliary data transmitted in accordance with the application of the present invention are referred to herein as robust VSB data (RVSB)."}
{"text": "The present invention relates to a novel hard, highly abrasion and corrosion-resistant, hydrogenated carbon material useful as an ultra-thin protective overcoat layer for high areal recording density thin film magnetic and magneto-optical (xe2x80x9cMOxe2x80x9d) recording media, a method of depositing films or layers of the novel material, and to improved magnetic and MO media including a protective overcoat layer comprised of the novel material.\nA magnetic recording medium, e.g., a hard disk, typically comprises a laminate of several layers, comprising a non-magnetic substrate, such as of Alxe2x80x94Mg alloy or a glass or glass-ceramic composite material, and formed sequentially on each side thereof, a polycrystalline underlayer, typically of chromium (Cr) or Cr-baed alloy, a polycrystalline magnetic recording medium layer, e.g., of a cobalt (Co)-based alloy, a hard, abrasion-resistant, protective overcoat layer, typically containing carbon (C), and a lubricant topcoat.\nIn operation of the magnetic recording medium, the polycrystalline magnetic recording medium layer is locally magnetized by a write transducer, or write head, to record and store information. The write transducer creates a highly concentrated magnetic field which alternates direction based on the bits of information being stored. When the local magnetic field produced by the write transducer is greater than the coercivity of the recording medium layer, then the grains of the polycrystalline recording medium at that location are magnetized. The grains retain their magnetization after the magnetic field produced by the write transducer is removed. The direction of the magnetization matches the direction of the applied magnetic field. The magnetization of the polycrystalline recording medium can subsequently produce an electrical response in a read transducer, allowing the stored information to be read.\nThin film magnetic recording media are conventionally employed in disk form for use with disk drives for storing large amounts of data in magnetizable form. Typically, one or more disks are rotated on a central axis in combination with data transducer heads. In operation, a typical contact start/stop (CSS) method commences when the head begins to slide against the surface of the disk as the disk begins to rotate. Upon reaching a predetermined high rotational speed, the head floats in air at a predetermined distance from the surface of the disk due to dynamic pressure effects caused by air flow generated between the sliding surface of the head and the disk. During reading and recording operations, the transducer head is maintained at a controlled distance from the recording surface, supported on a bearing of air as the disk rotates, such that the head can be freely moved in both the circumferential and radial directions, allowing data to be recorded on and retrieved from the disk at a desired position. Upon terminating operation of the disk drive, the rotational speed of the disk decreases and the head again begins to slide against the surface of the disk and eventually stops in contact with and pressing against the disk. Thus, the transducer head contacts the recording surface whenever the disk is stationary, accelerated from the static position, and during deceleration just prior to completely stopping. Each time the head and disk assembly is driven, the sliding surface of the head repeats the cyclic sequence consisting of stopping, sliding against the surface of the disk, floating in the air, sliding against the surface of the disk, and stopping.\nAs a consequence of the above-described cyclic CSS-type operation, the surface of the disk or medium surface wears off due to the sliding contact if it has insufficient abrasion resistance or lubrication quality, resulting in breakage or damage if the medium surface wears off to a great extent, whereby operation of the disk drive for performing reading and reproducing operations becomes impossible. The protective overcoat layer is formed on the surface of the polycrystalline magnetic recording medium layer so as to protect the latter from friction and like effects due to the above-described sliding action of the magnetic head. Abrasion-resistant, carbon (C)-containing protective coatings have been utilized for this purpose, and are typically formed by sputtering of a carbon target in an argon (Ar) atmosphere. Such amorphous carbon (a-C)-containing protective overcoat layers formed by sputtering have relatively strong graphitic-type bonding, and therefore exhibit a low coefficient of friction in atmospheres containing water (H2O) vapor, which characteristic is peculiar to graphite. However, the a-C layers produced in such manner have very low hardness as compared with many ceramic materials such as are employed as slider materials of thin film heads, and thus are likely to suffer from wear due to contact therewith.\nIn recent years, therefore, carbon-based protective overcoat layers having diamond-like hardness properties (i.e., HV of about 1,000-5,000 kg/mm2) have been developed, and films of diamond-like carbon (DLC) having a high percentage of diamond-type Cxe2x80x94C bonding have been utilized. Such DLC films exhibit a high degree of hardness due to their diamond-like sp3 bonding structure, and in addition, exhibit the excellent sliding properties characteristic of carbon, thus affording improved sliding resistance against sliders composed of high hardness materials. Such DLC films are generally obtained by DC or RF magnetron sputtering of a carbon target in a gas atmosphere comprising a mixture of Ar gas and a hydrocarbon gas, e.g., methane, or hydrogen gas. The thus-obtained films exhibit DLC properties when a fixed amount of hydrogen is incorporated therein. Incorporation of excessive amounts of hydrogen in the films leads to gradual softening, and thus the hydrogen content of the films must be carefully regulated.\nAmorphous, hydrogenated carbon films (referred to herein as a-C:H films) obtained by sputtering of carbon targets in an Ar+H2 gas mixture exhibiting diamond-like properties have also been developed for improving the tribological performance of disk drives; however, the electrical insulating properties of such type films lead to undesirable electrical charge build-up or accumulation during hard disk operation which can result in contamination, glide noise, etc. In order to solve this problem without sacrifice or diminution of the advantageous mechanical properties of such a-C:H films, attempts have been made to dope or otherwise incorporate nitrogen (N) atoms into the a-C:H films, in view of a substantial decrease in electrical resistivity and optical band gap (EBG) exhibited by such nitrogen-doped a-C:H films relative to undoped films.\nHowever, the continuous increase in areal recording density of magnetic recording media requires a commensurately lower flying height. Therefore, it would be advantageous to reduce the thickness of the carbon-based protective overcoat layer without adverse consequences. Conventional sputtered a-C:H materials are difficult to uniformly deposit and generally do not function satisfactorily at a thickness of about 30 xc3x85 or less. Specifically, conventional sputtered a-C:H films of about 30 xc3x85 thickness fail to provide adequate protection against corrosion of the underlying magnetic layer(s), particularly Co-containing ferromagnetic layers, when under environments of high temperature and humidity, and the resulting corrosion product(s) frequently are disadvantageously transferred to the transducer heads, often leading to failure of the disk drive.\nThe use of alternative deposition techniques for developing thinner and harder a-C:H layers having the requisite mechanical and tribological properties has been studied, such as chemical vapor deposition (CVD), ion beam deposition (IBD), and cathodic arc deposition (CAD) techniques. For example, the IBD method can be utilized for forming hydrogenated ion-beam carbon films (referred to herein as i-C:H films) that exhibit superior tribological performance at thicknesses below about 100 xc3x85. However, such films are insulating and, thus, suffer from the above-described drawback of electrical charge build-up during hard disk operation associated with sputtered a-C:H films.\nAccordingly, there exists a need for an improved hard, abrasion and corrosion-resistant material particularly suitable for use as an ultra-thin protective overcoat layer in high areal density magnetic recording media, and a method for manufacturing same, which method is simple, cost-effective, and fully compatible with the productivity and throughput requirements of automated manufacturing technology.\nThe present invention fully addresses and solves the above-described problems attendant upon the formation of ultra-thin, abrasion and corrosion-resistant protective overcoat layers suitable for use with high areal density magnetic recording media, such as are employed in hard drive applications, while maintaining full compatibility with all mechanical and electrical aspects of conventional disk drive technology. In addition, the present invention enjoys utility in the formation of ultra-thin, abrasion and corrosion-resistant protective overcoat layers required in the manufacture and use of thin film-based, ultra-high recording density magneto-optical (MO) data/information storage and retrieval media in disk form and utilizing conventional Winchester disk drive technology with laser/optical-based read/write transducers operating at flying heights on the order of a few micro-inches above the media surface.\nAn advantage of the present invention is an improved hard, abrasion and corrosion-resistant, hydrogenated carbon (xe2x80x9cC:Hxe2x80x9d) material formed by a simultaneous sputter+plasma enhanced chemical vapor deposition (xe2x80x9cPECVDxe2x80x9d) process.\nAnother advantage of the present invention is an improved protective overcoat material for magnetic and MO recording media and comprising an improved hard, abrasion and corrosion-resistant, sputter+PECVD-deposited C:H material.\nYet another advantage of the present invention is an improved magnetic or MO recording medium including an ultra-thin protective overcoat layer comprised of a sputter+PECVD-deposited C:H material.\nStill another advantage of the present invention is an improved simultaneous sputter+PECVD method for forming C:H films or layers suitable for use as ultra-thin, abrasion and corrosion-resistant protective overcoat materials in magnetic and MO recording media applications.\nA further advantage of the present invention is an improved method for regulating or controlling the amount of PECVD-derived carbon (C) atoms contained in sputter+PECVD deposited C:H materials, films or layers.\nA still further advantage of the present invention is an improved apparatus for performing simultaneous sputtering+PECVD for forming films or layers of the improved C:H materials according to the present invention.\nStill another advantage of the present invention is an improved magnetic recording medium comprising at least one ferromagnetic layer containing Co and means for protecting the at least one Co-containing ferromagnetic layer from corrosion under high temperature, high humidity conditions.\nAdditional advantages and other features of the present invention will be set forth in the description which follows and in part will become apparent to those having ordinary skill in the art upon examination of the following or may be learned from the practice of the present invention. The advantages of the present invention may be realized and obtained as particularly pointed out in the appended claims.\nAccording to one aspect of the present invention, the foregoing and other advantages are obtained in part by a novel hard, abrasion and corrosion-resistant material useful as an ultra-thin protective overcoat layer for a magnetic or magneto-optical (MO) recording medium, which novel material comprises hydrogenated carbon (C:H) formed by a process comprising simultaneous sputter and plasma-enhanced chemical vapor deposition (PECVD) of the hydrogenated carbon (C:H), wherein the amount of carbon atoms in the novel C:H material derived from the PECVD component of the process is less than about 50 at. %.\nAccording to certain embodiments of the present invention, the amount of carbon atoms in the novel C:H material derived from the PECVD component of the process is at least about 30 at. %; the position of the Raman G-band of the C:H material is about 1553 cmxe2x88x921; and the film resistance of the C:H material is as high as about 85 k xcexa9.\nAnother aspect of the present invention is a magnetic or MO recording medium comprising a protective overcoat layer formed of the novel C:H material; and according particular embodiments of the present invention, the thickness of the protective overcoat layer is not greater than about 30 xc3x85.\nYet another aspect of the present invention is a method of forming a layer of a novel hard, abrasion, and corrosion-resistant hydrogenated carbon (C:H) material on a surface of a substrate, which method comprises the steps of:\n(a) providing a vacuum chamber including a carbon sputtering target in the interior space thereof;\n(b) providing a substrate in the interior space of the chamber, such that a surface of the substrate is in facing relation to the sputtering target;\n(c) supplying the interior space of the vacuum chamber with at least one hydrocarbon gas and at least one inert gas at separately controllable flow rates and applying a sufficient negative potential to the carbon sputtering target to generate a plasma in the interior space to deposit a layer of a novel hard, abrasion and corrosion-resistant C:H material on the substrate surface by simultaneous sputtering of the carbon sputtering target and plasma enhanced chemical vapor deposition (PECVD) of carbon and hydrogen from the at least one hydrocarbon gas, wherein:\nstep (c) includes separately controlling the flow rates of each of the hydrocarbon and inert gases supplied to the vacuum chamber such that the amount of carbon atoms in the layer of C:H which are derived from PECVD of the hydrocarbon gas is less than about 50 at. %.\nIn accordance with embodiments of the present invention, step (c) includes separately controlling the flow rates of each of the hydrocarbon and inert gases to the vacuum chamber such that the amount of carbon atoms in the layer of C:H which are derived from PECVD of the at least one hydrocarbon gas is at least about 30 at. %; step (c) further includes supplying the interior space of the vacuum chamber with at least one hydrocarbon gas of formula CxHy, where x=an integer from 1 to 5 and y=an integer from 2 to 10 and at least one inert gas selected from the group consisting of He, Ne, Ar, Kr, and Xe.\nAccording to particular embodiments of the present invention, step (c) includes supplying the interior space of the vacuum chamber with at least one of acetylene (C2H2) and ethylene (C2H4) as the at least one hydrocarbon gas and Ar as the at least one inert gas; step (a) comprises providing an elongated, cylindrical carbon sputtering target rotatable about its axis of elongation; step (c) further comprises rotating the cylindrical carbon sputtering target about the axis of elongation; and step (b) further comprises applying a bias voltage within the voltage range from 0 to about xe2x88x92150 V to the substrate during step (c).\nEmbodiments of the present invention include comprises providing a magnetic or magneto-optical (MO) recording medium as the substrate in step (b), the surface thereof comprising the exposed surface of an uppermost layer of a stack of layers comprising the medium; and step (c) comprises forming a protective overcoat layer of the hard, abrasion and corrosion-resistant C:H on the exposed surface of the uppermost layer of the medium.\nAccording to particular embodiments of the present invention, step (b) comprises providing a disk-shaped substrate; and step (c) comprises forming the protective overcoat layer to a thickness not greater than about 30 xc3x85.\nStill another aspect of the present invention is a recording medium, comprising:\n(a) a substrate;\n(b) a stack of thin film layers on the substrate; and\n(c) a protective overcoat layer on an uppermost layer of the stack of thin film layers, the protective overcoat layer comprising a novel hard, abrasion and corrosion-resistant material comprised of hydrogenated carbon (C:H) formed by a process comprising simultaneous sputter and plasma-enhanced chemical vapor (PECVD) deposition of the novel hydrogenated carbon (C:H) material, wherein the amount of carbon atoms in the C:H material derived from the PECVD component of the process is less than about 50 at. %.\nAccording to embodiments of the present invention, the amount of carbon atoms in the C:H material derived from the PECVD component of the process is at least about 30 at. %; the position of the Raman G-band of the C:H material of the protective overcoat layer is about 1553 cmxe2x88x921; and the film resistance is as high as about 85 k xcexa9.\nIn accordance with particular embodiments of the present invention, the stack (b) of thin film layers comprises a stack of layers for a magnetic or magneto-optical (MO) recording medium; and the substrate (a) is disk-shaped.\nAccording to specific embodiments of the invention, the stack (b) of thin film layers comprises a stack of layers for a magnetic recording medium; the protective overcoat layer (c) is not greater than about 30 xc3x85 thick; and the stack (b) of thin film layers includes at least one ferromagnetic layer comprising Co.\nA further aspect of the present invention is an apparatus for performing simultaneous sputter and plasma-enhanced chemical vapor (PECVD) deposition of a layer of a novel hard, abrasion and corrosion-resistant, hydrogenated carbon (C:H) material on a surface of a substrate, comprising:\n(a) a vacuum chamber defining an interior space;\n(b) a carbon sputtering target in the interior space of the vacuum chamber, the carbon target being in the form of an elongated cylinder rotatable about its axis of elongation;\n(c) substrate mounting means for mounting a substrate in the interior space of the vacuum chamber, such that a major surface of the substrate is in parallel, facing relation to the elongated cylinder sputtering target; and\n(d) gas supply means for supplying the interior space of the vacuum chamber with at least one hydrocarbon gas and at least one inert sputtering gas at separately controllable flow rates.\nAccording to embodiments of the present invention, substrate mounting means (c) comprises means for mounting a disk-shaped substrate such that a major surface thereof is in parallel, facing relation to the cylindrical sputtering target; and the apparatus further comprises a bias voltage applying means (e) for applying to the substrate a bias voltage within the range from 0 to about xe2x88x92150 V.\nA still further aspect of the present invention is a magnetic recording medium, comprising:\n(a) at least one ferromagnetic thin film layer containing Co; and\n(b) means for protecting the at least one Co-containing ferromagnetic thin film layer from corrosion under high temperature, high humidity environments.\nAdditional advantages and aspects of the present invention will become readily apparent to those skilled in the art from the following detailed description, wherein embodiments of the present invention are shown and described, simply by way of illustration of the best mode contemplated for practicing the present invention. As will be described, the present invention is capable of other and different embodiments, and its several details are susceptible of modification in various obvious respects. Accordingly, the drawings and description are to be regarded as illustrative in nature, and not as limitative."}
{"text": "To achieve an increase in circuit density of integrated circuits, the size of semiconductor devices, such as field-effect transistors, within such integrated circuits has decreased. Decreasing the size of a semiconductor device can, however, result in a reduction in the length of a channel of the semiconductor device. Reducing the channel length can result in a source region and a drain region of the semiconductor device being closer to one another, which can allow the source and drain region to exert undue influence over the channel, or rather over carriers within the channel, commonly referred to as short-channel effects. Consequently, a gate of a semiconductor device that suffers from short-channel effects has reduced control over the channel, which, among other things, inhibits the ability of the gate to control on and/or off states of the semiconductor device."}
{"text": "The present invention relates to constructing a MONOWELL and more particularly to apparatus and methods for constructing a monodiameter wellbore for a monodiameter casing and a monobore production delivery system and still more particularly to drilling and completing a well using apparatus and methods to achieve a monodiameter wellbore, installing monodiameter casing and liners and installing a fullbore production delivery system.\nTraditional well construction, such as the drilling of an oil or gas well, includes a wellbore or borehole being drilled through a series of formations. Each formation, through which the well passes, must be sealed so as to avoid an undesirable passage of formation fluids, gases or materials out of the formation and into the borehole or from the borehole into the formation. In addition, it is commonly desired to isolate both producing and non-producing formations from each other so as to avoid contaminating one formation with the fluids from another formation.\nAs the well is drilled deeper, conventional well architecture includes casing the borehole to isolate or seal each formation. The formation may also be cased for borehole stability due to the geo-mechanics of the formation such as compaction forces, seismic forces and tectonic forces. The casings prevent the collapse of the borehole wall and prevent the undesired outflow of drilling fluids into the formation or the inflow of fluids from the formation into the borehole. The borehole may need to be cased due to equivalent circulating density and hydraulics reaching or exceeding the formation pore pressure or exceeding the fracture gradient pressure thus allowing fluids or gases to transfer between formations and borehole. If the formations are non-producing, or not of the desired producing interval, (some intervals are producing but at low levels) the formations can be cased together. If shallow water flows (where water flows several hundred feet below the seabed floor), or if there is potential communication among formations, then the formation is cased. The casings extend downhole and are sequentially placed across the formations through which the wellbore or borehole passes. The casings may be liners which do not extend to the surface of the wellbore. Traditionally steel casing has been used to case off formations.\nIn standard practice, each succeeding casing placed in the wellbore has an outside diameter significantly reduced in size when compared to the casing previously installed, particularly to accommodate hangers for the inner strings, and may be described as a series of nested casing strings. The borehole is drilled in intervals whereby a casing, which is to be installed in a lower borehole interval, is lowered through a previously installed casing of an upper borehole interval. As a consequence of this procedure, the casing of the lower interval has a smaller diameter than the casing of the upper interval. Thus, the casings are in a nested arrangement with casing diameters decreasing in the downward direction.\nThe use of a series of casings, which have sequentially reduced diameters is derived from long experience. The number of casings required to reach a given target depth is determined principally by the properties of the formations penetrated and by the pressures of the fluids contained in the formations. If the driller encounters an extended series of high pressure/low pressure intervals, the number of liners required under such circumstances may be such that the well cannot usefully be completed because of the continued reduction of the casing diameters required. Along with the downsize serial casing operations, the production tubulars may have to be downsized as well further reducing the delivery capacity of the well.\nIf the borehole extends through a formation that tends to cave in and thus causes the borehole to be very unstable, casing inserts must be installed to keep the borehole open. A casing insert is a type of emergency casing string which shores up an unstable formation and is an additional section of casing that is set through this unstable portion of the borehole. By requiring a casing insert for this unstable formation, an even smaller size casing than was planned is then required to complete the well. This reduces the diameter of the well and thus the ultimate internal diameter available for the production tubulars. The casing insert may not be possible, requiring that the well be sidetracked, resulting in a substantial reduced diameter wellbore.\nThe disadvantages of nesting casing and liners is apparent in slim-hole drilling. A slim-hole well is one in which 90% or more of the length of the well is drilled with bits smaller than 7 inches in diameter. See SPE 19525: An Innovative Approach to Exploration and Exploitation Drilling: The Slim-Hole High-Speed Drilling System by Walker and Millheim, September 1990. Slim hole drilling focuses on starting with a small borehole and finishing with an even smaller borehole for production.\nThe casing is fixed in the borehole by a cement layer between the outer wall of the casing and the wall of the borehole. During the drilling of the wellbore, annuli are provided between the outer surfaces of the casings and the borehole wall and a composition, sometimes referred to as “oil field” cement, is introduced in the annulus for cementing the casing within the wellbore. The casing is commonly cemented in place after the installation of each casing. When the casing is located in its desired position in the well, a cement slurry is pumped via the interior of the casing and around the lower end of the casing and upwards into the annulus, thereby causing the cement slurry to drive the drilling fluid upward in the annulus. As soon as the annulus around the casing is sufficiently filled with the cement slurry, injection of cement into the well is stopped and the cement slurry is allowed to harden. The cement sets up in the annulus, supporting and positioning the casing and forming a substantially impermeable barrier which divides the well bore into subterranean zones.\nUltimately the borehole reaches the target and is drilled through a hydrocarbon-containing formation or reservoir to produce hydrocarbons. The borehole may not be cased through the hydrocarbon-containing reservoir to allow substantially unrestricted influx of fluids from the formation into the borehole. When the formation is so weak that it will collapse, the uncased borehole section is completed with a liner. It is common practice to install a liner in the reservoir by suspending the liner in the borehole through the reservoir and then pumping a cement slurry into the annulus. After the cement has set to a hardened mass, perforations are extended through the liner and the cement body into the hydrocarbon-containing formation around the well in order to allow in-flow of reservoir hydrocarbon fluids, such as oil or gas, into the well.\nThe liner may be provided with slots to allow fluid influx into the borehole. The liner is usually secured at its upper end to the lower end of the non-productive barrier casing previously installed in the borehole. Because the slotted liner must pass through the previously set casing, it must have an outer diameter which is less than the inner diameter of the cased section. Over time, the formation may collapse and settle against the outer wall of the liner so that the area around the liner gets filled with particulates. U.S. Pat. Nos. 5,366,012 and 5,667,011 teach an expandable liner which is expanded by an expansion mandrel by moving the mandrel through the liner to radially expand the liner to a larger diameter in the borehole.\nThe purpose of the cement body around the casing is to fix the casing in the well and to seal the borehole around the casing in order to prevent vertical flow of fluid alongside the casing towards other formation layers or even to the earth's surface. Casing is traditionally cemented in place for two main reasons. (i) one to seal off and prevent leak paths between permeable zones and/or surface, and (ii) to give support and stability to the casings. The cement prevents fluid exchange between or among formation layers through which the wellbore passes, and prevents the undesirable migration of fluids between zones or gas from rising up the wellbore. It is important that there is no gas or fluid leakage after the cement has set and the well is completed.\nA problem generally encountered during cementation of the casing is, that due to various factors, such as the existence of varying pressure and temperature gradients along the length of the casing and shrinkage of the cement body during hardening thereof, relative displacements occur between the casing and the hardened cement mass which may result in poor bonding or cracking between the cement body and the casing. Poor bonding may result in the presence of a so-called micro-annuli between the casing and cement body, micro-annuli may extend along a substantial part of the length of the casing. The occurrence of a micro-annuli is particularly dangerous in gas wells as substantial amounts of gas might escape to the surface. In some cases hydrogen sulfide or natural gas can escape into the atmosphere. This condition may also lead to surface or ground water contamination. The resulting problems are very expensive to correct.\nThe poor bonding of the cement may be attributed to drilling fluid contamination or to bonding of the cement to the casing after the cement has set and/or oil or mill finish contamination on the surface of the casing or it can be attributed to aggressive drilling or aggressive pressure subjection and large pressure differs prior to it hardening and during the operation. As is well known in the art, hardening of cement causes generally a slight reduction of the volume of the cement. A more fundamental cause is the loss of hydrostatic head during the curing of the cement such that the formation pressure exceeds the annulus pressure and gas migration occurs causing channeling of the cement and subsequent leakage. Various additives and application techniques relative to the cement have been used in order to reduce the occurrence of this problem. During cementing operations, it is common to both reciprocate and rotate the casing during the cement pumping operation to break up or close any cement channels around the casing. Also compressible cement slurries have additives that entrain gas, which during the cement pumping operation, are compressed and as the hydrostatic head is lost during curing of the cement, the entrained gas subsequently expands and prevents loss of the pore pressure such that formation gas is prevented from migrating into the annulus. This technique, however, results in a lower strength cement. Thixotropic cement slurries depend on the cement achieving high gel strengths in very short time periods. If there is a rapid static gel strength obtained, gas migration and channeling are reduced or prevented. These specialized cement additives are expensive and require specific operational techniques. Thus, it is essential that a good bonding be created between the cement body and both the casing and the borehole wall.\nThere are various types of wells such as land based wells and offshore wells. Well applies to anything that produces oil, gas, water, or hydrates. Offshore wells may be shallow or deep water wells. A shallow offshore well is typically drilled from a platform that is in water up to 3,000 feet in depth. A deep water well is drilled from a floating platform or vessel with a riser extending from the sea floor to the platform or submersible rig. Any water deeper than 5,000 feet requires a drilling vessel, typically a drill ship.\nVarious types of casing may be installed in the well including conductor casing, surface casing, intermediate or production casing and production liners. Typically a land based well starts with a 20″/18⅝″ or larger diameter casing and telescopes down through two or three intermediate casings, to a final casing size of typically 6⅜″ with a 5″ production liner installed. Each casing is secured in place with cement filling an annulus having a size typically varying from 1 to 10 inches over the length of the casing and may be as much as 14 to 21 inches or greater at a wash out in the borehole wall.\nFIG. 1 is a schematic of a conventional deep water well completion. The size and number of casing and tubing strings will increase or decrease depending upon the well plan based upon, for example, the depth of the well, the production tubing delivery size, the structural support and the seabed formation support. If the seabed formation is unconsolidated and has little support, then the structural or conductor casing is larger and is set deeper. If the initial conductor casing is in rock, then it can be smaller with substantially less depth. For example, initially a structural or conductor casing and riser are lowered from a drilling platform and driven, drilled or jetted into the sea floor to provide support for a surface casing. The structural or conductor casing may or may not be cemented.\nFIG. 1 illustrates a 36 inch by 16 inch by 10¾ inch by 7 inch casing program with the addition of one or more tubing strings. After the 36 inch conductor casing is set, one or more surface casings is installed. A borehole is drilled for a 20 inch surface casing which is lowered into place with a 21″ surface casing riser attached thereto. A subsea wellhead with blowout prevention equipment, such as an 18¾ inch blowout preventer, is installed on the surface casing. The subsea wellhead may be supported by a structural casing.\nFurther, a borehole may be drilled through the riser and wellhead and through a problematic formation to extend a structural casing through the problem formation. For example, there are salt formations in the deepwater of the Gulf of Mexico. The structural casing forms a barrier across the formation while also supporting the wellhead. The structural casing has a thicker wall and provides a stable support frame for and can carry the load on the subsea wellhead. A 16 inch structural casing may be drilled, installed and cemented through a salt formation to seal off the salt formation from the wellbore being drilled. It should be appreciated that if there is no problematic formation, such as a salt zone, a shallow water flow zone, loss circulation zone, or other problem zone, then a structural casing is not needed to seal off the problematic area but it will support the subsea or platform wellhead, depending on well type.\nAnother borehole is then drilled for a 13⅜ inch intermediate casing string which is lowered into the borehole, attached to another riser, and cemented in place. Next a borehole may be drilled for another intermediate casing, such as a 11⅞ inch casing, and cemented in place. The borehole for the production casing string, such as a 9⅝ inch casing, is drilled and the production string is landed. It may or may not be cemented in place. The drilling is performed through blowout prevention equipment.\nSubsequently, a production tubing is installed and is supported within the wellhead on a tubing hanger, a hanger system or anchor system. The production tubing is typically 3½ inch tubing but may be as small as 1½ inches or as large as 12 inches. After the tubing hanger seals have been tested, the blowout prevention equipment is removed and a Christmas tree or subsea tree is installed. If the well is land based or being drilled from a platform, the blowout prevention equipment is at the surface. If the well is a subsea well, the blowout prevention equipment and tree are installed subsea. Also, the tubing can be installed through the subsea tree or subsea template. Thus, conventional techniques use a plurality of concentric casing strings with varying diameter and do not have a monodiameter architecture. It should be appreciated that conventional well architectures may vary depending upon the geological or drilling conditions.\nAs a consequence of the nested arrangement of the casings, a relatively large diameter borehole is required at the upper part of the borehole. Because the upper casing(s) has to be larger than the lower casing(s) for the lower casing(s) to pass through the upper casing(s), the upper portion of the borehole typically has a much larger diameter than the intended ultimate diameter at the bottom of the borehole. Large boreholes are disadvantageous in that they generate large amounts of cuttings and require increased volumes of drilling fluid and cement. In the standard well casing configuration, large volumes of cuttings are produced initially and heavy logistics are required during early phases of drilling. Generally speaking, larger borehole sizes take longer to drill than smaller diameter boreholes at equivalent depth. For example, increased drilling rig time is involved due to required cement pumping and cement hardening. Further, a large borehole diameter often takes larger fluid and horsepower capacity rigs generating increased costs due to heavy casing handling equipment and large drill bits. Thus conventional equipment results in larger boreholes drilled for each formation, larger sized equipment, greater fluid volumes, and larger casing strings than is absolutely required to provide a borehole for a well, for an injecting or producing or monitoring.\nUtilizing a large borehole often causes the usage of a wide variety of equipment and fluids that might not achieve maximum efficiency for the drilled borehole. If problems arise, additional fluids must be pumped and additional cement must be used to cement the formation to overcome the variances encountered during conventional well construction, otherwise a side track must be performed.\nConventional well architecture, engineering, and planning accounts for potential problem migration, well plan variance, and contingency. Therefore, large tolerances in equipment and procedures are provided in anticipation of variances in the length and/or composition of the formations, geomechanics, and growth/loading design. Compensation in the well architecture, engineering, and planning must be included in the well plan for contingency due to such large tolerances, including drilling for additional casing strings for geomechanical problems and sidetracks and re-drills for the installation of casing inserts prior to reaching the reservoir formation.\nIt can be appreciated that the problems with conventional well architecture are exacerbated in a deepwater well. In addition to the larger boreholes to be drilled, the larger sized completion equipment, the greater fluid volumes, and the larger casing strings, a deepwater well also requires large risers extending to the water's surface. The risers require the use of additional large fluid volumes, such as for drilling fluids and cement, to drill and cement the casing strings. Further, the large risers add substantial expense and additional large sized completion equipment.\nIt has long been an objective to achieve a monodiameter well where the wellbore is drilled from spud to total depth using one borehole size. For example the monodiameter well might be spudded with a driven conductor 7⅝″ to 9⅝″ in diameter. Thereafter, a borehole is drilled for each borehole section, perhaps a 7″ diameter. The borehole is then cased or lined with expandable casing or liners. Cement or some other innovative sealing composition is then used for the annular pressure seal. The next section of borehole is drilled using the same sized overgauge hole drilling and then cased or lined again with the same size expandable casing or liner. The process is repeated to target depth. See SPE 65184: “Towards a Mono-Diameter Well—Advances in Expanding Tubing Technology” by Benzie, Burge, and Dobson presented at the SPE European Program Committee Conference held Oct. 24-25, 2000.\nThe monodiameter well is designed based on the borehole size required across the reservoir. Rig capacity and all the drilling and completion equipment for the entire well is sized to the reservoir borehole size. Upon the advent of the monodiameter well, the telescoping well design with all its associated and myriad selection of drilling and completions equipment will become obsolete. The monodiameter well will achieve dramatic reductions in well construction costs. The challenge to the industry is to develop the full suite of enabling and complimentary technologies that will be required to drill and complete a monodiameter well. This suite of equipment will include drilling equipment, reaming while drilling (RWD), bi-center bits, energy balanced bits, near bit reamers, open hole annular sealing, well control procedures, well control equipment, and wellheads among others.\nExpandable tubulars are being developed whereby casings and liners are expanded diametrically after they are placed in wellbores. The ultimate use of expanded tubulars is in a monodiameter well, whereby the entire well is drilled and cased using effectively one hole size. A solid steel tubular can be readily expanded using forces, either mechanical or hydraulic, available on most drilling and workover rigs. Expandable tubulars can be used in open hole either as a temporary drilling liner or as a permanent liner tied back to the previous casing string. See SPE 54508: “The Reeled Monodiameter Well” by Pointing, Betts, Bijleveld, and Al-Rawahi, presented at the 1999 SPE/CoTA Coiled Tubing Roundtable, May 25-26, 1999, and SPE 65184: “Towards a Mono-Diameter Well—Advances in Expanding Tubing Technology” by Benzie, Burge, and Dobson, presented at the SPE European Petroleum Conference, Oct. 24-25, 2000.\nThe concept of expanding a solid tube is relatively simple. A mandrel or “pig” whose outside diameter is greater than the tube's inside diameter is forced through the tube, thereby plastically deforming the tubing material to a larger diameter size. A solid tubular can be expanded using a cone of ceramic, tungsten carbide or hardened steel that is either mechanically pulled through the tubular or hydraulically pushed. The solid tubular can be expanded to around 30%-40%, although a range from 10%-20% will probably be more typical. The combination of expansion cone radius, material characteristics, expansion ratio, annular seal material, and gauge tolerance of the open borehole or casing within which the tubular is to be expanded, all determine the expansion forces and the tolerances or fit of the final expanded tubular in the wellbore. During the expansion process, the tubular strength increases since the expansion process is a cold working of the material. However, the collapse strength of a post expanded tubular is less than that of the pre-expanded tubular but it is within the design limits expected for the required borehole pressures.\nOne of the requirements of expandable casing/liners is to run a drill bit through the casing and drill a hole of a greater diameter than the previous casing. The next section of casing will then be run through the previous one and expanded against it. Thus, an enlarged wellbore having a diameter greater than the external diameter of the preceding installed casing must be drilled such as by “overgauge hole drilling” which encompasses the use of expandable bits, bi-center or eccentric bits or equivalent, reaming while drilling (RWD), under-reamers or similar tools and other novel drilling methods known to those skilled in the art, such as expandable/retractable stabilizers. Typical bits for overgauge hole drilling include bi-center bits and eccentric bits. A bi-center bit has a well-defined pilot bit section and an eccentric wing mounted further back on the bit body. An eccentric bit resembles more of a conventional bit, but a flank on the high-lobe side has a longer profile than on the other side. Eccentric bits are generally used in softer formations whereas bi-center bits are more commonly used in harder formations. The major difference of both bi-center and eccentric bits over under-reamers is that a section of borehole can be drilled to the required size in one run. However, an under-reamer provides a smoother and more diametric borehole wall.\nAt a desired depth, or when it is otherwise decided to case or line and cement the wellbore, an expandable casing or liner, whose greatest external (outside) diameter approximates, i.e., is only slightly smaller than the internal diameter of the casing or liner previously installed, is lowered through the previously installed casings or liners and into the newly drilled open enlarged borehole. The lowered casing or liner is a tubular made of a deformable material. The deformable casing/liner may have a decreased wall thickness. The lowered casing/liner is positioned in relation to the wellbore so that the upper end of the lowered casing/liner overlaps the lower end of the previously installed casing.\nIn expanding the expandable casing/liner, a die member, such as a mandrel or cone, is drawn or pumped through the length of the lowered casing to expand the casing in-situ. The die member has a suitable shape and composition, such as hardened steel, and is adapted or sized and shaped to expand the liner to the diameter of the previously installed casing. The die member is shaped or designed to provide at least substantially uniform expanded or deformed liner segment of circular or approximately circular periphery.\nThe upper end of the die member is connected to a running string and is pulled through the casing. Further, the die member may have a fluid tight seal, such as a cupseal, for sealing the die member within the casing and allowing sufficient fluid pressure to produce movement of the die member. Any suitable wellbore fluid or liquid available can be used for displacing the die member. To ease the expansion process, the inside diameter of the casing may be lubricated to allow the die member or expansion cone to move smoothly through the casing. The rate of upward adjustment or movement of the die member by upward movement of the running string and the application of pressure below the die member may be correlated so as to reduce movement of the die member up through the casing with concurrent gradual deformation and expansion of the casing, providing an expansion achieving an external diameter equal to or approximating, preferably slightly greater or larger than, that of the previously installed casing.\nThreaded connections of lengths of expandable casing or liner remain the primary connection of choice. However, the connection has to be internally flush to allow the die member to pass through the connection, and externally flush to allow expansion to occur with constant expansion force. SPE 54508: “The Reeled Monodiameter Well” by Pointing, Betts, Bijleveld, and Al-Rawahi, presented at the 1999 SPE/CoTA Coiled Tubing Roundtable, May 25-26, 1999, discloses using coiled casing that maintains a single diameter of well bore throughout. Coiled casing can be used in the production of a monodiameter wellbore as well as tubular jointed expandable casing. Coiled casing can be expanded or installed non-expanded. A reeled monodiameter casing or liner has the same throughbore.\nThe expansion causes a virtual forced fit at the overlap of the upper end of the expanded lower casing/liner into the lower end of the previously installed casing. Thus, hangers are not required since the casing/liner is supported by the previously installed string. A constant through-bore is maintained at the overlap. The well is then subsequently completed by internally cladding the last casing string. Cladding the overlapped portions of the adjacent casings allows the lower casing to be supported at the cladding by the upper casing and allows the upper and lower casings to be pressure sealed at the overlap. The pipe over the area (overlap) to be cladded may include a corrosion resistant coating. One type of expandable tubular is disclosed in U.S. Pat. No. 6,085,838, hereby incorporated herein by reference.\nThe new well concept of using expandable casing and liners necessitates a narrow annuli in the range of 3 to 4 inches on diameter or less between the casing/liners and the borehole wall. Thus a quality borehole and optimum borehole size are required. The wellbore or borehole must have diametric efficiency. Diametric efficiency is the maintenance of the optimum hole size regardless of other well construction requirements or constraints. Preferably diametric efficiency maintains the optimum borehole size from the surface through the producing reservoir. The use of expandable casing/liners requires the maintenance of diametric efficiency to achieve a monodiameter well. Diametric solutions that maintain or improve diametric efficiency in a range of applications from multilaterals, high pressure/high temperature (HPHT), extended reach, horizontals and deepwater environments and remedial applications (i.e. side tracks) must be developed to use expandables.\nReferring now to FIG. 2, there is shown a bell-end design disclosed in the SPE 54508 article. The annulus is split up into two parts: (i) where the liner is expanded against the previous casing/liner, either full length or overlap, and (ii) where the casing is expanded over the open borehole section.\nThe annular clearances between the casing/liner and borehole wall must allow cement slurries to displace the drilling fluid effectively. The interface or overlap between casing/liner strings must have the mechanical ability of the interface to hold axial loads and must have the hydraulic ability to form a pressure tight seal between the two casing/liner strings.\nA simple expansion of metal against metal does not provide a reliable seal. Elastomeric seals are used to provide a pressure seal between the overlapped casing ends. Inside a previous casing or a gunbarrel hole through hard formations, a polymer or an elastic rubber coating may be applied externally to the casing/liner. This can deform elastically with the expansion of the casing/liner and form a seal between the casing/liner and/or formation.\nA thicker layer of heat softening rubber, which can be coated on the outside of the expandable pipe, may be used in slightly overgauge boreholes. The thermal plastic rubber is reformed in-situ to its final shape by performing the casing/liner expansion process at a more elevated temperature, achievable by heating the well either before or during the expansion process.\nPreviously, a variety of cement compositions have been used for cementing in open hole applications. However, cement is undesirable for use with expandable casing. Oilfield cement or hydraulic cement compositions are incompressible and tend to resist the expansion of the casing or liner making the expansion more difficult. Thus, expansion of the casing/liner can lead to the crushing of the cement, and consequent loss of effectiveness regarding the zones. This problem is exacerbated by the small annular clearance associated with expandable casing/liners. With the new well concepts as described and the narrow annuli created in the range of 3 to 4 inches on diameter or less, other isolation materials must now be considered. Isolation materials must be ductile in order to form an appropriate seal. Conventional cement will be particularly brittle and weak in such thin sheaths or thicknesses and is hence an inappropriate sealing medium. Therefore, a composition with comparable strength to cement, but with greater elasticity and compressibility is required for cementing expandable casing.\nOther problems may be encountered using oilfield cement or hydraulic cement as the sealing composition for expandable casing/liners. If the cement composition gels or sets prior to accomplishing the expansion, the cement composition is crushed in the annular space between the walls of the well and the expandable casing or liner whereby it does not function to seal the expanded casing or liner in the wellbore.\nA conventional cement sealant composition under expansion conditions, with appropriate cement recipe to prevent the pre-maturely setting of the cement, may be used. Sealant compositions are required for sealing expandable casings or liners in well bores. Such compositions are compressible and maintain the properties required to provide a seal between the walls of the wellbore and the expanded casings or liners. Several types of sealants materials have been tested. Two materials are silicon gel and a two-component silicon rubber. Suitable wetting agents may be used such as barite or iron micropellets. The silicon gel is a highly viscous material. The two-component silicon rubber consists of two materials mixed in equal volume and weight. The mixture will automatically set, but a retarder may be added to control setting times. Silicon gels or rubbers are very ductile.\nPrior to the expansion, the sealant material may be pumped into the annulus between the walls of the well bore and the unexpanded casing or liner in a similar manner to conventional cementing techniques. Once in place, the casing is expanded and the material allowed to set. In a severely washed out borehole, such as soft formations with wash outs and ledges, the sealant material may not be adequate.\nIn the traditional well, the wellbore may be open 3 to 5 days before stability of the borehole wall becomes a problem. After that, portions of the borehole wall may start to fall into the borehole. In using expandable casing/liners, longer intervals are going to be drilled thus leaving the borehole open for longer than 3 to 5 days. Thus there is a much longer time period than in conventional drilling during which the borehole wall is unsupported by casing during drilling. It is important that the borehole wall remain stable for these longer periods of time. Therefore, special drilling fluids must be used to drill the enlarged boreholes for expandable casing/liners to ensure the stability of the borehole over longer time periods.\nDrilling fluid is circulated downwardly through the drill string, through the drill bit and upwardly in the annulus between the walls of the well bore and the drill string. The drilling fluid functions to remove cuttings from the well bore and to form a filter cake on the borehole wall. As the drilling fluid is circulated, a filter cake of solids from the drilling fluid forms on the walls of the well bore. The filter cake build-up is a result of initial fluid loss into permeable formations and zones penetrated by the well bore. The presence of the filter cake reduces additional fluid loss as the well is drilled.\nIn addition to removing cuttings from the well bore and forming filter cake on the well bore, the drilling fluid cools and lubricates the drill bit and maintains a hydrostatic pressure against the well bore walls to prevent blow-outs, i.e., to prevent pressurized formation fluids from flowing into the well bore when formations containing the pressurized fluids are penetrated. The hydrostatic pressure created by the drilling fluid in the well bore may fracture low mechanical strength formations penetrated by the well bore which allows drilling fluid to be lost into the formations. When this occurs, the drilling of the well bore must be stopped and remedial steps taken to seal the fractures. Such remedial actions are time consuming and expensive.\nIt is preferred to drill as long an interval as possible without having to stop drilling and case the borehole. However, in order to insure that fracturing of low mechanical strength formations penetrated by the well bore and other similar problems do not occur, it may become necessary to case and cement the borehole. As previously discussed, it is preferred to avoid frequently installing and cementing casing because drilling must be stopped and frequent casing of the borehole may cause a reduction in the producing borehole diameter.\nAnother problem that occurs in the drilling and completion of well bores is that when the well bore is drilled into and through unconsolidated weak zones or formations. Unstable materials such as clays, shales, sand stone and the like make up a high percentage of the formations in which wells are drilled, and a majority of well bore problems are a result of the instability of such materials, particularly shale instability. Shales are sedimentary rocks that contain a variety of clays. Shales containing montmorillonite, often referred to as smectite clays, swell and disperse when contacted by water. Shales which swell upon contacting water are often referred to as heaving or sloughing shales. Such shales upon contact with aqueous drilling fluids swell and fracture rendering the well bore wall unstable. In such cases, the well bore wall sloughs into the well bore. Sloughing of shale and other similar unstable materials into the well bore can cause the drill string to become stuck and can enlarge the well bore resulting in large subterranean cavities. Additionally, when sloughing occurs while the drill bit is being changed at the surface, the well bore fills up and must be cleared before drilling can proceed.\nFurthermore, the heaving unstable material suspended in the drilling fluid increases its solid content, and as a result, the viscosity of the drilling fluid increases to the point where the drilling fluid must be chemically treated to reduce its viscosity or it must be diluted followed by the addition of weighting material to maintain its mud weight. The instability of clays, shales, sand stone and the like is also caused by hydraulic pressure differential leading to fluid transport and by pressure changes near the well bore as the drilling fluid compresses pore fluid and diffuses a pressure front into the formation.\nConsolidating unconsolidated weak zones or formations formed of clays, shales, sand stone and the like while drilling a well bore prevents sloughing of the clays, shales, sand stone and the like into the well bore and prevents the need for implementing time consuming and costly remedial steps. It is preferred to increase the mechanical strength of the well bore whereby hydrostatic pressure exerted on the well bore by the drilling fluid does not cause fractures or the like to occur in the well bore. Such fractures cause drilling fluid to be lost and also require stoppage of the drilling operation and costly remedial steps to be taken.\nAnother significant advantage of increasing the mechanical strength of the well bore is the reduction or elimination of casing intervals at which casing or liners are cemented in the well bore which reduces or eliminates the overall time and cost of cementing the well. An additional advantage is that the well bore has a larger diameter in the production zone, due to fewer casing intervals, which increases productivity.\nA monobore is a term used in the industry for a fullbore production delivery system. OTC 8585: “Case History, Well Completion and Servicing Strategies for the Hibernia Field” by Wylie, Maier, Shamloo, Huffman, and Downton presented at the 1998 Offshore Technology Conference, May 4-7, 1998, defines a monobore. To address overall completion criteria, a monobore design was chosen for the Hibernia Field for the initial completions. This design allowed flowbore access across the well's pay zone without diameter restrictions (not necessarily constant diameter).\nIn conventional completions, there are internal valves and gauges which project into the bore of the completion tubular. These narrow the inside diameter of the completion tubulars. Each restriction in the completion tubulars causes a pressure drop across that point causing scale and corrosion to occur at these points. The narrowing of the completion tubulars also reduces flow and therefore production.\nThe present invention overcomes the deficiencies of the prior art"}
{"text": "In recent years, technologies that deal with materials having nanometer size, namely nanomaterials, have been attracting public attention in various industrial fields. New materials having unprecedented new functions have been developed by combining the nanomaterials with a plurality of other materials at the nanometer level. Since nanomaterials dispersed in high degree show different properties from those in a bulk state and hence they are useful, technology to disperse nanomaterials in a composite is indispensable to obtain the composite in which nanomaterials are dispersed in high degree. However, there is a problem that peculiar functions of nanomaterials cannot be realized since they agglomerate when combined to make composite materials because the surface state of the nanomaterials is generally unstable.\nAmong the nanomaterials, a carbon nanotube has been studied for its physical properties and functions since its discovery in 1991, and research and development for its application have also been carried out extensively. However, in the case of combining the carbon nanotube with a resin or a solution to make a composite material, there is a problem that its characteristic properties cannot be realized because it further agglomerates from its entangled state when it has been produced. Therefore, various attempts to uniformly disperse or dissolve the carbon nanotube into a solvent or a resin have been tried by carrying out physical treatment or chemical modification on the carbon nanotube. For example, a method of dispersing single-walled carbon nanotube by cutting it shortly through its ultrasonic wave treatment in a strong acid has been proposed (Non-patent Document 1). However, this is not an industrially proper method because the operation is complicated owing to the treatment in the strong acid, and the effect of this method on dispersion is insufficient.\nSuch a shortly cut single-walled carbon nanotube has open edges, each terminated with an oxygen-containing functional group such as a carboxylic group, and taking account of this, it has been proposed to make the shortly cut single-walled carbon nanotube solubilizable in a solvent by changing the carboxylic group into an acid chloride and reacting the acid chloride with an amine compound to introduce a long chain alkyl group into the shortly cut single-walled carbon nanotube (Non-patent Document 2). However, this method has problems unsolved in which there occurs a damage of a graphene sheet structure of the carbon nanotube and in which characteristic properties of the carbon nanotube itself are affected because the long chain alkyl group is introduced into the single-walled carbon nanotube by a covalent bond.\nAs another attempt, a method has been reported in which, using the fact that pyrene molecule is adsorbed on the surface of the carbon nanotube through a strong mutual interaction, a substituent group containing an ammonium ion is introduced into the pyrene molecule and the resultant molecule is subjected to ultrasonic treatment in water together with the single-walled carbon nanotube so that the resultant molecule is non-covalently adsorbed on the single-walled carbon nanotube to thereby produce a water soluble single-walled carbon nanotube (Non-patent Document 3). According to this method, the damage of a graphene sheet structure and the like can be suppressed owing to a noncovalent bonding type chemical modification, however, there is a problem that electrical conductivity of the carbon nanotube is lowered caused by the existence of a nonconductive pyrene compound.\nA method has been proposed in which a dispersion liquid is obtained by dispersing or solubilizing the carbon nanotube in a solvent such as water or an organic solvent without deteriorating the characteristic properties of the carbon nanotube itself using a general-purpose surfactant or polymer dispersant (Patent Document 1 and Patent Document 2). There is a description that the carbon nanotube is stably dispersed in the dispersion liquid, however, there is no description about a dispersed state of the carbon nanotube in a coating film or a composite to be formed from the dispersion liquid or about an application to a conductive material or the like.\nFurther, there has been proposed a composition composed of carbon nanotube, a conductive polymer, and a solvent, and a composite produced from it (Patent Document 3). It is reported that the composition and the composite can disperse or solubilize the carbon nanotube in a solvent such as water, an organic solvent, or a water-containing organic solvent and can be excellent in long-term storage stability without deteriorating the characteristic properties of the carbon nanotube itself by the coexistence of the conductive polymer. The carbon nanotube composition in which the conductive polymer coexists is excellent in electrical conductivity, film-formability, and moldability and can be applied or coated on a base material with an easy method, however, application to materials that need colorlessness and transparency is difficult owing to a coloring originated from the conductive polymer to be used. Further, there has been a problem that usable solvents are limited because the conductive polymer generally has low solubilities to various solvents.\nFurther, there has been proposed a composition composed of a nanomaterial, a (meth)acrylic polymer, and a solvent, and a composite produced from it (Patent Document 4). It is reported that the composition and the composite can disperse or solubilize the nanomaterial in a solvent such as water, an organic solvent, or a water-containing organic solvent and can be excellent in long-term storage stability without deteriorating the characteristic properties of the nanomaterial itself by the coexistence of the (meth)acrylic polymer. The nanomaterial composition in which the (meth)acrylic polymer coexists is excellent in film-formability and moldability and can be applied or coated on a base material with an easy method, however, there is a problem that an usable solvent is limited because an amine compound has to be jointly used to raise dispersibility. Further, in the case of using the nanomaterial composition as a coating film (cured film) by adding it to a polymerizable monomer, there has been a problem that water resistance and mar resistance as a coating film are lowered since the nanomaterial composition has no crosslinking point with the film component because the nanomaterial composition is a compound containing sulfonic group and the like.\nFurther, there has been proposed a photosensitive composition for a black matrix, composed of a binder resin containing a urethane acrylate compound obtained by a reaction between an isocyanate compound and a polyhydroxy compound and a black pigment such as carbon black or carbon nanotube (Patent Document 5). It is described that the photosensitive composition containing the urethane acrylate compound can produce the black matrix at a low cost because it can easily form a pattern having a high light blocking effect with a thin film and has a high sensitivity and a high curing speed, however, there is no description about a dispersion effect of the carbon nanotube and an additional dispersant is separately added to disperse the black pigment. Further, the black matrix has to realize the light blocking effect by a coating film obtained from a composition containing a black pigment, and hence, in the case of applying the composition to a use where transparency is required, a dispersion level of the carbon nanotube is low and accordingly a coating film having a sufficient transparency cannot be obtained. Further, there is a problem that nothing is described concerning an effect of improving electrical conductivity obtained by sufficiently dispersing the carbon nanotube.    Patent Document 1: International Publication No. WO 2002/016257    Patent Document 2: Japanese Patent Application Laid-Open No. 2005-35810    Patent Document 3: International Publication No. WO 2004/039893    Patent Document 4: International Publication No. WO 2006/028200    Patent Document 5: Japanese Patent Application Laid-Open No. 2005-331938    Non-patent Document 1: R. E. Smalley, et. al., Science, 280, 1253 (1998)    Non-patent Document 2: J. Chen, et. al., Science, 282, 95 (1998)    Non-patent Document 3: Nakajima, et. al., Chem. Lett., 638 (2002)"}
{"text": "The present invention is a brick block and process and apparatus for its manufacture. This invention is particularly useful for manufacturing chimney block having a brick pattern of preferably three courses impressed on all four vertical sides as well as for manufacturing wall block which may have a brick pattern on one, two or three sides.\nThe brick block of this invention comprises a rectangular block of monolithic concrete colored throughout to imitate the color of brick. At least one vertical face of the rectangular block is horizontally and vertically grooved to create a pattern simulating laid brick. The grout lines contain a coloring agent (e.g. a pigmented latex) that simulates the color of a mortared joint. The brick block may be a chimney block, all four vertical sides of which have the impressed brick pattern. It may also be a wall block having two cores and having one to two sides faces with the brick pattern or if it be an end wall block having three sides of the block with the pattern.\nBy the term xe2x80x9cbrick blockxe2x80x9d is meant a molded monolithic concrete block colored throughout to simulate a natural brick color one vertical face of which has an impressed brick grout pattern simulating laid brick, one, two, and preferably at least three or more courses high.\nThe apparatus for the manufacture of the brick block comprises a mold having four vertical walls the interior of which has the dimensions of the desired brick block. One or more of the walls has means for impressing the grout pattern on the face of the brick block as it is molded. This means includes a hydraulic system that moves a grout bar outwardly into the concrete mass as it is being formed approximately xe2x85x9 inch or so and after it is formed and before the block is ejected from the mold retracts the grout bar to be flush with the mold face.\nThe block has a core as is usual in the manufacture of conventional wall and chimney block. The concrete mixture has zero slump so that the impressed grout pattern is retained in the block as it is ejected from the mold. The concrete mixture has an increased moisture content while maintaining zero slump to produce on the xe2x80x9cbrickxe2x80x9d face of the block, a brick-like surface. This results in the block having a higher density as compared to that of conventional wall block. The density is usually above 100 pounds per cubic foot, e.g. 103 pounds per cubic foot.\nEjection is accomplished by having a head or ram descend from the top of the mold around the core to push the block downwardly onto a platen which is then moved out from underneath the mold and set off for curing of the concrete mass.\nThe hydraulic mechanism that extends and retracts the grout bar has two or more hydraulic cylinders attached to the bar so that the grout bar remains square with the face of the mold.\nThe mold of the present invention is adapted to fit in or onto a conventional block manufacturing machine such as the xe2x80x9cBescopacxe2x80x9d manufactured by Besser, 81 Johnson Street, Alpena, Mich. 49707. Because of the hydraulic mechanism required to impress the grout pattern, the mold is oversized or larger than that used for a standard block and some slight modification has to be made to the machine to receive the frame holding the present mold, namely, the width of the rubber dam is decreased.\nOnce the mold is in place, manufacture is quite straightforward. A charge of concrete, pre-colored throughout, is pushed into the mold cavity around a core or cores held in place by overhead hangers. As is conventional, the concrete is caused to settle in place by vibration during which time the grout bar is extended into the cement mixture and then is retracted after vibration has ceased. Following this, a head or ram descends into the mold to push the block out of the mold onto the platen on which it is removed. The core may be vented to prevent any vacuum forming which would distort the block. After the block has been allowed to cure, the grout lines of the brick pattern are painted as with a masonry slur or a pigmented latex the color of mortar.\nThe cycle time per brick block is 12 to 15 seconds. During a six-hour operating period, 1100 brick block have been made using one operator."}
{"text": "1. Field\nEmbodiments of the invention relate to electronic devices, and more particularly, to notch filters.\n2. Description of the Related Technology\nAn amplifier, such as an operational amplifier, can include chopper circuitry and a notch filter to help reduce the amplifier's input offset voltage. For example, the amplifier can include a first chopping circuit used to periodically swap or chop the positive and negative inputs of the amplifier's input stage at a chopping frequency and a second chopping circuit used to periodically swap or chop the positive and negative outputs of the amplifier's input stage at the chopping frequency. Additionally, a notch filter can be used to generate a notched signal by attenuating frequency components of the second chopping circuit's output signal that are at the chopping frequency. Thereafter, the notched signal can be amplified to generate the amplifier's output signal.\nThe second chopping circuit can operate to upshift the input offset voltage of the amplifier's input stage from DC to the chopping frequency where it can be attenuated by the notch filter. Since the frequency spectrum of the input signal can be upshifted by the first chopping circuit and downshifted by the second chopping circuit, the input signal can pass through the notch filter without attenuation.\nAlthough configuring the amplifier to include the chopper circuitry and the notch filter can reduce the amplifier's input offset voltage, mismatch error in the notch filter can result in the amplifier nevertheless having offset and/or the amplifier's output signal having ripples at the chopping frequency and at harmonics thereof.\nThere is a need for notch filters having improved performance. Additionally, there is need for chopper-stabilized amplifiers having reduced input offset voltage and/or reduced output voltage ripple."}
{"text": "The present invention relates to a method and apparatus for controlling at a base station the data transmission and retransmission of a mobile station, and in particular to a method and apparatus for, in a wireless communication network, controlling new data transmission and retransmission of a mobile station by a base station with minimal signaling overhead in a downlink and performing communication without any data-associated signaling in an uplink.\nThe 3GPP (3rd Generation Partnership Project) launched LTE (Long Term Evolution) in 2005, which aims to provide support for increasing requirements from operators and subscribers with higher data throughput and better network performance.\nThe 3GPP R6 has introduced High Speed Downlink Packet Access (HSDPA) and High Speed Uplink Packet Access (HSUPA). The primary design target for a wireless communication system is packet field service to support highly real-time service and reduce delay in wireless network. The current criterion is derived with reference mainly to Quality of Services (QoS) of various real-time services, such as Voice over IP (VoIP) and the like. Meanwhile, status transition delay of a control plane in the system is reduced, which is helpful to improve users' experience of network services. The air unidirectional delay is less than 5 ms in the user data plane, and the delay for a user transiting from an idle status to a connection status is less than 100 ms in the control plane. The 3GPP evolution involves technical standards of CDMA2000, WCDMA (Wideband Code Division Multiple Accesses) and LTE.\nFIG. 1 shows schematically a mobile communication network by example of a LTE system. As shown in FIG. 1, in the mobile communication network, a plurality of base stations (NodeBs) are connected with a core network gateway (MME/UPE). One core network gateway manages several NodeB nodes. Each of the plurality of base stations can adapt more than one mobile station (i.e., User Equipment, UE), and each mobile station can always perform transmission/reception to/from the respective base station by using a dedicated channel (Dedicated Physical Control Channel, DPCCH). Further, the mobile stations in the mobile communication network can perform data reception with HSDPA and data transmission with EUDCH (Enhanced Uplink Dedicated Channel). The base stations conduct schedule based on these packet transmission systems, and allocate radio resources to the respective mobile stations for communication with the mobile stations.\nIn order to achieve an error-free data transmission, the error control strategy adopted in LTE is still Hybrid ARQ (HARQ) obtained from a combination of FEC (Forward Error Coding) and ARQ. HARQ utilizes a retransmission strategy of Increment Redundancy (IR). Retransmission serves as an approach to error correction in the data link layer of mobile communication, and HARQ has a good performance in this respect.\nTo improve system performance, HSUPA adopts the technique HARQ at the physical layer and supports two types of combining mechanisms, i.e., Chase Combining (CC) of the same packets retransmitted from a base station and IR Combining of packets containing different information (i.e., redundant information) retransmitted from a base station. Information is transmitted directly between UE and NodeB in a fashion of ACK/NACK, in which the receiving side will send ACK information to the transmitting side over a respective channel if it receives data correctly, otherwise, it will send NACK information. This facilitates the transmitting side to learn whether a retransmission is needed in an accurate and timely way.\nIn current phase, LTE uplink transmission has the following characters:\n1. The 3GPP RAN 1 work group has determined not to use any data-associated signaling in uplink transmission;\n2. Synchronized HARQ is adopted in LTE uplink transmission;\n3. It has been proposed now that a common Maximum HARQ retransmission number is applied for all radio bearers;\n4. The 3GPP has now specified that adaptive HARQ can be used for uplink transmission; and\n5. The downlink signaling overhead should be minimized to increase the system capacity.\nIn HSUPA, the data-associated signaling, RAN=0, is used in the uplink to indicate to eNodeB (base station) that the mobile station is now performing new data transmission, which aims to ensure a correct HARQ combination by the eNodeB, so that the base station can determine, based on the quality of received signals, whether the mobile station will conduct new data transmission or data retransmission in the next step. In LTE, however, the RAN 1 work group has determined not to use any data-associated signaling in uplink transmission. As a result, the mobile station (UE) cannot notify the eNodeB whether it is conducting a new data transmission or data retransmission. It is thus impossible to ensure a correct HARQ combination by the eNodeB.\nIn HSUPA, the data-associated signaling is used in the uplink to indicate a new data transmission, for example, notifying the receiving side how the data are transmitted via signaling. The receiving side can demodulate the data using a corresponding method. Unfortunately, the prior art strategy cannot be applied in LTE since the use of any data-associated signaling is not allowed in LTE uplink. Occasions may occur that the receiving side (base station) is not aware that the transmitting side (UE) has begun to transmit new data and still waits for retransmission of previous data from the transmitting side. To keep synchronization between the transmitting and receiving sides in a certain fashion, the receiving side shall learn whether the transmitting side is transmitting new data or retransmitting data."}
{"text": "X-ray detectors for use in medical computerized tomography are well known in the art. Medical x-ray detectors, such as described and claimed in commonly assigned U.S. Pat. No. 4,031,396 to Whetten et al., are subject to the resolution-limiting effects of k-band x-ray fluorescence. This phenomenon relates to the emission of x-ray photons from the incident x-ray energy-excited gas molecules of the detecting medium and is fully explained in the above-cited patent. The atoms of a gaseous detecting medium, upon interacting with the incident x-ray energy, may themselves emit low energy x-ray photons. When the detecting medium allows a long absorption length for a fluorescent x-ray photon, the photon thus generated by secondary emission may be detected within the detector at a location remote from the point where the incident x-ray entered the detector. Such event degrades spatial resolution and decreases detector efficiency since a portion of the incident x-ray energy is thus not usefully detected. To address this problem, the Whetten et al. patent discloses a detector embodiment which employs voltage and collector plates comprising metals of high atomic number, such as tungsten. Such plates are substantially opaque to x-ray energy and effective to absorb the fluorescent x-rays, thereby minimizing their adverse effect on resolution.\nThe use of tungsten plates, however, causes a decrease in operating efficiency of an x-ray detector. The detector operating efficiency may be expressed as the percentage of incident x-ray energy that is sensed by the detector. It will be appreciated that absorption of fluorescent x-rays by tungsten plates aborts the liberation of charged particles in the detecting medium which can be detected at the collector plates. Hence, that portion of incident x-ray energy responsible for the emission of fluorescent x-rays goes undetected.\nAnother drawback to using tungsten plates is that they can absorb incident x-ray energy just as well as fluorescent x-rays. In the above-cited Whetten et al. patent, the tungsten plates are oriented parallel to the incident x-ray beam. While these plates are quite thin (0.05 mm), x-ray energy striking a plate edgewise will be absorbed. Therefore, in the aggregate, the sum of the thicknesses of the plurality of voltage and collector plates used in the Whetten et al. detector represents a significant incident x-ray absorption area and materially contributes to loss of detector efficiency.\nTypical prior art medical x-ray detectors, such as disclosed in the Whetten et al. patent, employ a relatively low pressure gaseous detecting medium having a correspondingly low density. As a result, the absorption lengths of both the incident x-rays and fluorescent x-rays are relatively long. Such long absorption lengths necessitate the use of tungsten plates which extend a considerable distance into the gaseous detecting medium, e.g. approximatey one inch. The resulting large area presenting by each plate may cause a number of problems. First, microphonics are produced by any motion of the voltage plates with respect to adjacent collector plates. This motion may occur as a result of detector rotation in computerized axial tomography scanning, or even as a result of vibrations transmitted through the building structure supporting the equipment. Such voltage plate movements cause changes in the relative capacitances between the plates, which induce signals in the collector plates. The induced signals are superimposed on the true signals associated with the x-ray image and thus tend to distort spatial resolution. Moreover, the larger plate size must be accommodated, structurally, in order to maintain it rigid with respect to its adjacent collector plates.\nSecond, the large plates in the tungsten plate detector embodiment disclosed in the Whetten et al. patent are arranged in radial fashion aimed at the focal spot of the incident x-ray source. As a result, the position of the detector in relation to the x-ray source must be maintained precisely; otherwise, much of the incident beam will impinge on and be absorbed by the tungsten plates and hence not be detected.\nThird, use of larger voltage plates can result in degradation of kilovolt peak (KVP) linearity. KVP linearity relates in general to a detector's sensitivity to different energy spectra of incident x-rays. The individual collector elements of an x-ray detector are usually calibrated in accordance with a particular x-ray energy spectrum. A patient undergoing examination with a medical computerized tomography x-ray detector may be exposed to an incident x-ray energy spectrum comprising a range of energies. Absorption of x-rays by the patient's body results in a shifting of the average energy of the spectrum toward a higher value. In order to maintain acceptable tomographic image resolution, it is necessary that the individual collector elements respond uniformly to this shifting of the average energy. To this end, it is necessary that the volume of the detecting medium allocated to each collector element be exactly the same. When the respective portions of the detecting medium volume associated with the individual collector elements become unequal, KVP linearity is degraded. Uniformly increasing these volume portions improves KVP linearity. However, in prior art medical detectors, the large tungsten plates require extra care in manufacture to maintain, with precision, equality between the adjacent volume portions respectively associated with the collector elements.\nIndustrial x-ray detectors known in the art are typically employed in such applications as inspection of jet engine turbine buckets, nuclear fuel rods, weldments, and other industrial inspection functions. Typically, these applications require higher resolution and higher energy x-rays than do medical applications. For example, industrial detectors may be used at an average x-ray energy of 300 to 400 KEV while the average x-ray energy in a medical detector application is only on the other of 50 to 125 KEV.\nIn order to achieve the higher resolution required in industrial applications, such industrial detectors have utilized structures in which the individual detector elements are extremely narrow and closely spaced, e.g. individual detector widths of only 2 to 3 mils and spacing of only 1 to 2 mils. Such detector element structure provides a degree of spatial resolution that is neither necessary nor desirable in medical computerized tomography. If such a high resolution detector structure were to be applied in a medical x-ray detector, the lower signal-to-noise ratio that would be experienced would require an intolerable increase in the x-ray dosage to the patient under study in order to achieve acceptable results.\nIndustrial x-ray detectors typically employ a high density gaseous detecting medium in order to minimize the absorption length of the higher energy x-rays being detected. For example, the detector disclosed in U.S. Pat. No. 4,394,578 to Houston et al. employs a detecting medium gas pressure of up to 200 atmospheres.\nThe extremely high (1) detecting medium gas pressures, (2) image resolution structures and (3) x-ray energies are important industrial x-ray detector design considerations which are not applicable to medical x-ray detectors."}
{"text": "In order to make the urban and suburban environment more aesthetic but also more conducive to good healthy living, there is the need to introduce trees into the landscape.\nNonetheless it is well recognised that trees require a certain amount of space within these populated areas if they are going to develop into mature trees offering the benefits of shaded foliage, water retention, cooling, aesthetics and so forth to the surrounding area.\nFor the most part, town planning have seen the planting of trees in urban and suburban areas to grow in close proximity to pavements and walkways so that benefits of the trees can be enjoyed by those pedestrians utilising such places.\nNonetheless this is where the great conflict or competing demands begin to present themselves when trees are attempted to be placed in populated areas or positioned along pavements, walkways and so forth.\nIt is well recognised that the trees to successfully grow require rich nutrient soils with the appropriate levels of moisture and sufficient drainage to allow the tree's root system to pass there through.\nAlternatively walkways or pavements require compact well supported soil structure to keep the positioned pavers or concrete in place thereby avoiding any structural damage to the construction.\nHence if trees are going to be planted in urban and suburban areas there seems to be a predicament as to whether or not one needs to look after the tree and therefore provide low compact soils suitable for root growth or alternatively compact soils which provide the necessary load bearing support for the pavement and make it unconductive for the root system of the tree to develop and therefore extend out or damage the pavement in any way.\nIn order to try and overcome these kinds of problems more recently systems have been introduced that accommodate tree rooting, pavement and side walk developments in close proximity and also potentially offering further functionality for storm water management.\nFIG. 1 enclosed herewith provides such an example of how a modulated cell can create uncompacted soil volumes for large tree growths and also on site storm water management through the creation of a modular frame positioned about and below a load bearing feature such as a pavement and also within the tree root network.\nNonetheless problems with such an arrangement as the one depicted in FIG. 1 are quite numerous.\nAs the person skilled in the art will appreciate the rectangular crate type structure presented in FIG. 1 while it may have a lot of open spaces to accommodate a tree's root network, it does so at the expense of vertical stability of the frame structure when the cell frame is mounted vertically one upon the other.\nAlso the cell frame is made up of several independent components including the base, top plate and also the fastening means.\nWhile there is some scope of the cell frame to be mounted vertically, albeit quite efficiently due to the amount of space in between the beams providing support, there is no inherent lateral fastening mechanisms to bring together adjacent cell frames to form a matrix type configuration.\nStill further with so many individual components making up the cell frame it would be difficult to conveniently pack together a whole series of these frames during storage or transportation to their intended site for construction.\nWhile further efforts have been made by others to provide similar cell frames again to try and create some stackable frame structure to provide load bearing support but also accommodate a tree's root network, these other attempts have not been able to create a balance between the requirement of the load bearing aspects of supporting the pavement without detrimentally affecting the usefulness of the frame structure in allowing the root network of the tree to pass there through.\nWhile it may be possible to add further posts and beams and make the cell frame of reduced dimensions to the one illustrated in FIG. 1 this then takes away the benefit of open spaces and the like to allow the tree root network to pass there through and will also make the construction of the overall structural frame cumbersome, time consuming and expensive.\nTherefore there clearly remains a need in this area of interest associated with being able to provide a feature which is open to accommodate the needs of both a tree but also those physical engineering requirements of an appropriately level load bearing soil structure to support important transport sections such as side walks or walkways that must be free of any hazards such as cracks or tree root intrusions.\nAccordingly it is an object of this invention to provide a modulated cell which is adapted to form an improved structural frame which is able to provide support for a load bearing feature such as a pavement or walkway but at the same time adequately accommodate a rooting area for a tree.\nFurther objects and advantages of this invention including the ability of the same modulated cell structure to assist in storm water retention and the diversion of the storm water to underground tanks and so forth will become apparent from the complete reading of this specification."}
{"text": "This invention relates to an apparatus for controlling the drive of and/or the operation of a quenching car/bucket car for a coke oven as well as a method therefor, and more particularly to an apparatus for effectively accomplishing the automatic control of the operation of the quenching car/bucket car and a method therefor.\nConventionally, the control of the drive and/or operation of a quenching car for a coke oven is generally carried out by an operator riding in an electric car drawing the quenching car. Now, the manner or procedure of controlling operation of the quenching car will be described. First, when red-hot coke, as calcined is pushed out, by means of a pusher, from a designated one of the coking chambers of a coke oven from which the coke is to be delivered while keeping the pusher, a guide car and a quenching car accurately registered or aligned with respect to the designated coking chamber, it is charged onto the quenching car of which a door is kept closed while causing the quenching car to travel in correspondence with the speed of the coke pushed out. After the charging is completed, the quenching car is moved into a quenching tower and located accurately at a predetermined position in the quenching tower. Then, water is sprinkled from a water tank arranged above the quenching tower on the red-hot coke charged onto the quenching car to quench the coke and then the water is drained off from the coke. Subsequently, the quenching car is moved to coke wharves and then accurately stopped at a predetermined position in front of designated one of coke wharves. Thereafter, the door of the quenching car is opened to discharge the quenched coke onto the wharf. Thus, the processing of calcined coke is completed. The above-described procedure is repeated with respect to each of the coking chambers in order, resulting in the delivery of coke from the coking oven being carried out.\nThe control of drive of a bucket car for a coke oven is likewise carried out by an operator riding in an electric car drawing the bucket car. More particularly, red-hot coke as calcined is pushed out or discharged, by means of a pusher, from designated one of coking chambers of a coke oven from which coke is to be delivered while keeping the pusher, a guide car and a bucket car accurately located or aligned with respect to the designated coking chamber. The discharged coke is then charged into a coke bucket on the bucket car. Subsequently, the bucket car is moved to a coke dry quenching facility (hereinafter referred to as \"CDQ\") and accurately located at a predetermined position in CDQ. Then, the coke bucket in which the red-hot coke is charged is hung up by means of a charging crane to cause the coke to be subject to a dry quenching treatment. The coke bucket which has thus been rendered empty is then placed down into the bucket car for the next operation. The above-described operation is repeated with respect to each of the coking chambers, resulting in the delivery of coke from the coking oven to the bucket car being carried out.\nFor the operation of the quenching car/bucket car for the coke oven, the correlation between the quenching car/bucket car and other working machines such as the guide car, pusher and the like is highly important. To this end, it is generally required that workers or operators carefully and closely maintain contact with each other using a communication means such as a wire telephone, a wireless telephone or the like to verify with each other the working condition of the working equipment, such as its working position or the like.\nHowever, the artificial or manual operation of the quenching car/bucket car carried out while the operators maintain contact with each other is apt to cause exchange of incorrect information due to misunderstanding or the like, resulting in the operators often failing to carry out the operation while properly grasping the working situation, the working position and the like.\nIn view of the foregoing, an automatic control means is proposed for successively automatically controlling the operation of the working machines such as the pusher, guide cat, quenching car/bucket car and the like. However, taking, as an example, a track surface on which the working position of the quenching car/bucket car is determined, lots of load is applied onto the track surface because the sum of weight of the quenching car/bucket car and weight of the electric car traveling while drawing the quenching car/bucket car reaches a high level such as, for example, 90 tons (-55 tons (quenching car)+35 tons (electric car)). This causes the positional relationship of the track surface to a body of the coke oven, the quenching tower, the coke wharf/CDQ and the like to be readily varied, so that the track surface fails to serve as a reference surface for the operation of the quenching car/bucket car, resulting in a failure to permit the automatic control to be effectively carried out. Also, in order to permit the track surface to endure such high load as described above, it is required that there is a rigid concrete foundation to strengthen the track surface. Unfortunately, this leads to a problem of increased costs."}
{"text": "The present invention relates to viscoelastic surfactant fluids useful in subterranean operations, and more particularly, to dual-function additives that enhance fluid loss control and the stability of viscoelastic surfactant fluids, and their associated methods of use.\nTreatment fluids may be used in a variety of subterranean treatments, including, but not limited to, stimulation treatments and sand control treatments. As used herein, the term “treatment,” or “treating,” refers to any subterranean operation that uses a fluid in conjunction with a desired function and/or for a desired purpose. The terms “treatment,” and “treating,” as used herein, do not imply any particular action by the fluid or any particular component thereof. These subterranean operations include, but are not limited to, hydraulic fracturing treatments, acidizing treatments, gravel-packing treatments, sand control treatments, and the like.\nMaintaining sufficient viscosity in the treatment fluids used in these operations is important for a number of reasons. Maintaining sufficient viscosity is important in fracturing and sand control treatments for particulate transport and/or to create or enhance fracture width. Also, maintaining sufficient viscosity may be important to control and/or reduce fluid loss into the formation. At the same time, while maintaining sufficient viscosity of the treatment fluid often is desirable, it may also be desirable to maintain the viscosity of the treatment fluid in such a way that the viscosity also may be reduced easily at a particular time, inter alia, for subsequent recovery of the fluid from the formation.\nTo provide the desired viscosity, polymeric gelling agents commonly are added to the treatment fluids. The term “gelling agent” is defined herein to include any substance that is capable of increasing the viscosity of a fluid, for example, by forming a gel. Examples of commonly used polymeric gelling agents include, but are not limited to, guar gums and derivatives thereof, cellulose derivatives, biopolymers, and the like. To further increase the viscosity of a treatment fluid, often the polymeric gelling agent is crosslinked with the use of a crosslinking agent. Treatment fluids comprising crosslinked gelling agents also may exhibit elastic and/or viscoelastic properties, wherein the crosslinks between gelling agent molecules may be broken and reformed, allowing the viscosity of the fluid to vary with certain conditions such as temperature, pH, and the like.\nThe use of polymeric gelling agents, however, may be problematic. For instance, polymeric gelling agents may leave an undesirable gel residue in the subterranean formation after use, which can impact permeability of the formation. As a result, costly remedial operations may be required to clean up the fracture face and proppant pack. Foamed treatment fluids and emulsion-based treatment fluids have been employed to minimize residual damage, but increased expense and complexity often have resulted.\nTo combat perceived problems associated with polymeric gelling agents, some surfactants have been used as gelling agents. It is well understood that, when mixed with a fluid in a concentration above the critical micelle concentration, the molecules (or ions) of surfactants may associate to form micelles. The term “micelle” is defined to include any structure that minimizes the contact between the lyophobic (“solvent-repelling”) portion of a surfactant molecule and the solvent, for example, by aggregating the surfactant molecules into structures such as spheres, cylinders, or sheets, wherein the lyophobic portions are on the interior of the aggregate structure and the lyophilic (“solvent-attracting”) portions are on the exterior of the structure. These micelles may function, among other purposes, to stabilize emulsions, break emulsions, stabilize a foam, change the wettability of a surface, solubilize certain materials, and/or reduce surface tension.\nWhen used as a gelling agent, the molecules (or ions) of the surfactants associate to form micelles of a certain micellar structure (e.g., rodlike, wormlike, vesicles, etc., which are referred to herein as “viscosifying micelles”) that, under certain conditions (e.g., concentration, ionic strength of the fluid, etc.) are capable of, inter alia, imparting increased viscosity to a particular fluid and/or forming a gel. Certain viscosifying micelles may impart increased viscosity to a fluid such that the fluid exhibits viscoelastic behavior (e.g., shear thinning properties) due, at least in part, to the association of the surfactant molecules contained therein. As used herein, the term “viscoelastic surfactant” refers to surfactants that impart or are capable of imparting viscoelastic behavior to a fluid due, at least in part, to the association of surfactant molecules to form viscosifying micelles. Moreover, because the viscosifying micelles may be sensitive to hydrocarbons, the viscosity of these surfactant fluids may be reduced after introduction into the subterranean formation without the need for certain types of gel breakers (e.g., oxidizers). The term “breaker” is defined herein to include any substance that is capable of decreasing the viscosity of a fluid. This may allow a substantial portion of the surfactant fluids to be produced back from the formation without the need for expensive remedial treatments. Moreover, these viscoelastic surfactants may not leave the undesirable gel residue in the subterranean formation found in uses of polymeric gelling agents, reducing or alleviating the need for costly remedial operations.\nHowever, the use of viscoelastic surfactant fluids may be problematic in certain subterranean formations exhibiting high temperatures (e.g., above about 200° F.). Many viscoelastic surfactant fluids become unstable at these temperatures, which reduces the viscosity of the fluid. Moreover, the stability of viscosifying micelles in viscoelastic surfactant fluids may be extremely sensitive to various conditions (e.g., temperature, pH, presence of other additives in the fluid, composition of the subterranean formation, etc.), and thus the inclusion of other additives in the viscoelastic surfactant fluid that are needed for a given treatment using that fluid may detrimentally affect the rheological properties (e.g., viscosity) of the fluid. This inability to maintain a desired level of viscosity at higher temperatures, among other problems, may increase fluid loss and decrease the ability of the fluid to suspend and/or transport particulate materials.\nNumerous additives are used in the art to help control fluid loss in subterranean operations. Numerous additives are also used in the art to maintain stability and/or viscosity of a treatment fluid at higher temperatures. However, the use of these conventional additives may give rise to other problems. First, the necessity of both a fluid loss control additive and a separate stabilizing or viscosifying additive in a treatment fluid may increase the complexity and cost of a treatment fluid and/or a subterranean application utilizing that fluid. Moreover, many conventional fluid loss control additives permanently reduce the permeability of a subterranean formation, affect the rheology of the treatment fluid in which they are used, and/or reduce the rate at which the fluid is allowed to penetrate or leak off into the subterranean formation. However, in some instances, while it may be desirable to control or prevent fluid loss for a given period of time, it may become desirable to allow the treatment fluid to penetrate or leak off into the subterranean formation, or to increase the permeability of the subterranean formation, at some later point in time. Costly and time-consuming operations may be required to reverse the effects of conventional fluid loss control additives on the treatment fluid and/or to restore permeability to those portions of the subterranean formation affected by the fluid loss control additives."}
{"text": "Development of multicellular organisms depends, at least in part, on mechanisms which specify, direct or maintain positional information to pattern cells, tissues, or organs. Various secreted signaling molecules, such as members of the transforming growth factor-beta (TGF-β), Wnt, fibroblast growth factors and hedgehog families have been associated with patterning activity of different cells and structures in Drosophila as well as in vertebrates. Perrimon, Cell: 80: 517-520 (1995).\nHedgehog (Hh) was first identified as a segment-polarity gene by a genetic screen in Drosophila melanogaster, Nusslein-Volhard et al., Roux. Arch. Dev. Biol. 193: 267-282 (1984), that plays a wide variety of developmental functions. Perrimon, supra.; Hammerschmidt et al., Trends Genet. 13: 14-21 (1997). Although only one Drosophila Hh gene has been identified, three mammalian Hh homologues have been isolated: Sonic Hh (Shh), Desert Hh (DHh) and Indian Hh (IHh). Reviewed by Hammerschmidt et al., Trends Genet. 13: 14-21 (1997). Shh is expressed at high level in the notochord and floor plate of developing vertebrate embryos where it plays a key role in neural tube patterning. Echelard et al., Cell 75: 1417-30 (1993), Ericson et al., Cell 81: 747-56 (1995), Hynes et al., Neuron 19: 15-26 (1997), Krauss et al., Cell 75, 1431-44 (1993), Marti et al., Nature 375: 322-25 (1995), Roelink et al, Cell 81: 445-55 (1995). Shh also plays a role in the development of limbs (Laufer et al., Cell 79, 993-1003 (1994)), somites (Fan and Tessier-Lavigne, Cell 79, 1175-86 (1994); Johnson et al., Cell 79: 1165-73 (1994)), gut (Roberts et al., Development 121: 3163-74 (1995), lungs (Bellusci et al., Develop. 124: 53-63 (1997) and skin (Oro et al., Science 276: 817-21 (1997), as well as the regulation of left-right asymmetry (reviewed by Ramsdell and Yost, Trends in Genetics 14: 459-65 (1998)). Likewise, IHh and DHh are involved in bone and germinal cell development, Vortkamp et al., Science 273: 613-22 (1996), Bitgood et al., Curr. Biol. 6: 298-304. Shh knockout mice further strengthened the notion that Shh is critical to many aspect of vertebrate development, Chiang et al., Nature 383: 407-13 (1996). These mice show defects in midline structures such as the notochord and the floor plate, absence of ventral cell types in neural tube, absence of distal limb structures, cyclopia, and absence of the spinal column and most of the ribs.\nAt the cell surface, the Hh signals is thought to be relayed by the 12 transmembrane domain protein Patched (Ptch) [Hooper and Scott, Cell 59: 751-65 (1989); Nakano et al., Nature 341: 508-13 (1989)] and the G-protein coupled like receptor Smoothened (Smo) [Alcedo et al., Cell 86: 221-232 (1996); van den Heuvel and Ingham, Nature 382: 547-551 (1996)]. Both genetic and biochemical evidence support a receptor model where Ptch and Smo are part of a multicomponent receptor complex, Chen and Struhl, Cell 87: 553-63 (1996); Marigo et al., Nature 384: 176-9 (1996); Stone et al., Nature 384: 129-34 (1996). Upon binding of Hh to Ptch, the normal inhibitory effect of Ptch on Smo is relieved, allowing Smo to transduce the Hh signal across the plasma membrane. Loss of function mutations in the Ptch gene have been identified in patients with the basal cell nevus syndrome (BCNS), a hereditary disease characterized by multiple basal cell carcinomas (BCCs). Disfunctional Ptch gene mutations have also been associated with a large percentage of sporadic basal cell carcinoma tumors, Chidambaram et al., Cancer Research 56: 4599-601 (1996); Gailani et al., Nature Genet. 14: 78-81 (1996); Hahn et al., Cell 85: 841-51 (1996); Johnson et al., Science 272: 1668-71 (1996); Unden et al., Cancer Res. 56: 4562-5 (1996); Wicking et al., Am. J. Hum. Genet. 60: 21-6 (1997). Loss of Ptch function is thought to cause an uncontrolled Smo signaling in basal cell carcinoma. Similarly, activating Smo mutations have been identified in sporatic BCC tumors (Xie et al., Nature 391: 90-2 (1998)), emphasizing the role of Smo as the signaling subunit in the receptor complex for Shh.\nHowever, the exact mechanism by which Ptch controls Smo activity still has yet to be clarified and the signaling mechanisms by which the Hh signal is transmitted from the receptor to downstream targets is unclear. Genetic epistatic analysis in Drosophila has identified several segment-polarity genes which appear to function as components of the Hh signal transduction pathway, Ingham, Curr. Opin. Genent. Dev. 5: 492-98 (1995); Perrimon, supra.\nSignaling by hedgehog has been shown to be transduced in vertebrates through the Gli family of zinc finger transcription factors, Hynes et al., Neuron 19: 15-26 (1997); Lee et al., Development 124: 2537-52 (1997); Sasaki et al., Development 124: 1313-22 (1997); Ruiz, i Altaba, Development 125: 2203-12 (1998), and in Drosophila by the Gli homologue Cubitus interruptus (Ci) (Orenic et al., Genes Dev. 4: 1053-67 (1990); Alexandre et al., Genes Dev. 10: 2003-13 (1996); Dominquez et al., Science 272: 1621-25 (1996). Consistent with a pivotal role for Ci in transducing the Hh signal, several genes have been identified genetically in Drosophila and shown to modulate Ci activity (reviewed by Goodrich and Scott, Neuron 21: 1243-57 (1998); Ingham, Embo. J. 17: 3505-11 (1998). These include the putative serine threonine kinase fused (Fu), Preat et al., Genetics 135: 1047-62 (1993), a novel protein designated Suppressor of fused (Su(fu)) [Pham et al., Genetics 140: 587-98 (1995); Preat, Genetics 132: 725-36 (1992)] protein kinase A (PKA), Li et al., Cell 80: 553-562 (1995); Pan and Rubin, Cell 80: 543-52 (1995)], the kinesin-like molecule, Costal-2 (Cos-2) [Robbins et al., Cell 90: 225-34 (1997); Sisson et al., Cell 90: 235-45 (1997)], and the F-box/WD40 repeat protein slimb [Jiang and Struhl, Nature 391: 493-496 (1998)]. Additional elements implicated in Hh signaling include the transcription factor CBP [Akimaru et al., Nature 386: 735-738 (1997)], and the Shh response element COUP-TFII [Krishnan et al., Science 278: 1947-1950 (1997)].\nMutations in Cos-2 are embryonicly lethal and display a phenotype similar to Hh over expression, including duplications of the central component of each segment and expansion domain of Hh responsive genes. In contrast, mutant embryos for Ci of fused show a phenotype similar to Hh loss of function, while mutations in negative regulators of the Hh pathway, such as ptch or PKA, induce ectopic expression of Hh-target genes (reviewed by Ingham, Embo. J. 17: 3505-11 (1998)). For example, fused and Ci mutants exhibited deletion of the posterior part of each segment and replacement of a mirror-like image duplication of the anterior part or each segment and replacement of a mirror-like duplication of the anterior part, Busson et al., Roux. Arch. Dev. Biol. 197: 221-230 (1988). Molecular characterizations of Ci suggested that it is a transcription factor which directly activates Hh responsive genes such as Wingless and Dpp, Alexandre et al., (1996) supra, Dominguez et al., (1996) supra. Likewise, molecular analysis of fused reveals that it is structurally related to serine threonine kinases and that both intact N-terminal kinase domain and a C-terminal regulatory region are required for its proper function, Preat et al., Nature 347: 87-9 (1990); Robbins et al., (1997), supra; Therond et al., Proc. Natl. Acad. Sci. USA 93: 4224-8 (1996). However, whereas fused null mutations and N-terminal kinase domain mutations can be fully suppressed by Suppressor of fused mutations, C-terminus mutations of fused display a strong Cos-2 phenotype in a Suppressor of fused background. This suggests that the fused kinase domain can act as a constitutive activator of Shh signaling when Suppressor of Fused is not present.\nSu(fu) was originally isolated as a gene, which when activated, was able to suppress the embryonic and adult phenotypes of fused mutants, and when duplicated, enhanced the fused mutant phenotype, suggesting that fused and Su(fu) have antagonistic roles. [Preat, Genetics 132: 725-36 (1992); Preat et al., Genetics 135: 1047-62 (1993)]. Su(fu) mutant flies have a wing phenotype similar to but not as strong as patched or PKA mutants (Ohlmeyer and Kalderon, Nature 396: 749-53 (1998). The combination of patched or PKA mutations in a Su(fu) mutant background enhances the mutant phenotype of patched and PKA, suggesting a cooperative effect of these genes in modulating hedgehog signaling. Ohlmeyer and Kalderon, supra. Fused, Su(fu), Cos-2 and Ci have been shown to form a microtubule-associated multiprotein complex and hedgehog signaling leads to dissociation of this complex from microtubules. Robbins et al., Cell 90: 225-34 (1997); Sisson et al., Cell 90: 235-45 (1997); Monnier et al., Curr. Biol. 8: 583-86 (1998).\nBoth fused and Cos-2 become phosphorylated in response to Hh treatment, Robbins et al., supra; Therond et al., Genetics 142: 1181-98 (1996), but the kinase(s) responsible for this activity(ies) remain(s) to be characterized. To date, the only known vertebrate homologues for these components are members of the Gli protein family (e.g., Gli-1, Gli-2 and Gli-3). These are zinc finger putative transcription factors that are structurally related to Ci. Among these, Gli-1 was shown to be a candidate mediator of the Shh signal [Hynes et al., Neuron 15: 35-44 (1995), Lee et al., Development 124: 2537-52 (1997); Alexandre et al., Genes Dev. 10: 2003-13 (1996)] suggesting that the mechanism of gene activation in response to Hh may be conserved between Drosophila and vertebrates.\nIn the absence of hedgehog, full length Ci (Ci-155) is proteolytically processed into an N-terminal repressor fragment (Ci-75). Aza-Blanc et al., Cell 89: 1043-53 (1997). Recent studies demonstrate that complex formation is necessary to target Ci for proteolysis. Methot and Basler, Cell 96: 819-31 (1999). The cleavage of Ci potentially requires PKA phosphorylation of Ci and ubiquitination by Slimb, which targets Ci to the proteosome. Chen et al., Proc. Natl. Acad. Sci. USA 95: 2349-54 (1998); Jiang and Struhl, Nature 391: 493-96 (1998). In response to Hh, Ci cleavage is blocked and Ci-155 is activated into a labile but still uncharacterized form. Ohlmeyer and Kalderon, supra; Methot and Basler, Cell 96: 819-31 (1999).\nTo determine whether other signaling components in the Hh cascade are evolutionarily conserved and to examine the function of fused in the Hh signaling cascade on the biochemical level, Applicants have isolated and characterized human used cDNA, a kinase homologous the Drosophila Fu (dFu). Tissue distribution on the mouse indicates that fused is expressed in Shh and other hedgehog responsive tissues, and also displays the same subcellular localization as human Gli1 (hGli1) and hSu(fu), the human homologue of Drosophila Su(fu) (dSu(fu)).\nBiochemical studies demonstrate that fused is a functional kinase and that it forms a complex with hSu(fu) and hGli1. Functional studies provide evidence that fused is an activator of Gli and that a dominant negative form of fused is capable of blocking Shh signaling in Xenopus embryos. Applicant also herein show that Shh signaling leads to the reversible dissociation of human Su(fu) from human Gli-1 (hGli-1) in mammalian cells. Applicants also demonstrate herein that the catalytic subunit of protein kinase A (PKAc) is present in a complex in association with hSu(fu). PKAc phosphorylates both hSu(fu) and Gli, and thereby promotes the binding of hSu(fu) to Gli, while ectopic hFu or Shh stimulation trigger the dissociation of hSu(fu) from Gli. These biochemical observations correlate with data obtained in a functional readout where fused abrogates hSu(fu)-mediated repression of Gli in a Gli reporter assay. Together this data demonstrated generally that fused is directly involved in Hh signaling and specifically that fused antagonizes PKAc activity, thereby triggering the dissociation of hSu(fu) from hGli-1. Regulation of the hSu(fu)-hGli-1 interaction is central to the control of hGli-1 activity and is promoted by PKAc and inhibited by Shh and hFu.\nApplicants have identified a cDNA encoding a human fused (hfused) polypeptide and thus have provided for the first time a vertebrate fused molecule."}
{"text": "The purpose of a wiper pivot shaft air inlet grill seal is to prevent debris from entering the plenum and the heating, ventilating, and air conditioning (HVAC) system of the vehicle. The seal also prevents excessive water from falling on the windshield wiper motor, and provides a decorative cosmetic cap for the upper pivot housing. Typically, current vehicles are provided with a cosmetic cap for the pivot housing with no seal engaging the air inlet grill. An attempt is made to prevent infiltration of debris by making the aperture in the air inlet grill as small as tolerances will allow around the wiper pivot housing. Generally the aperture is 10 mm to 20 mm larger than the pivot housing diameter. Alternatively, an attempt is made to prevent infiltration of debris by adding a separate foam seal in addition to the decorative beauty cap placed on the pivot housing.\nIt would be desirable in the present invention to provide a decorative cap including an integral, one-piece, seal for engagement with the air inlet grill in order to prevent debris and fluid from entering the plenum and HVAC system of the vehicle. It should be desirable in the present invention to provide a snap-fit seal engagable with a windshield wiper pivot shaft and an air inlet grill to prevent excessive fluid from falling on the windshield wiper motor."}
{"text": "In addition to computers, home networks now typically include multiple types of subscriber equipment configured to deliver subscriber services through the home network. The subscriber services include the delivering of multimedia, such as streaming audio and video, through the home network to the subscriber equipment where it is presented to a user. As the number of available subscriber services increases, so does the number of devices being connected to a home network. The increase in the number of services and devices increases the complexity of the coordination between the network nodes as each node may be produced by a different manufacturer at different times. Some home networking technologies have emerged in an attempt to facilitate simple home network solutions and take advantage of existing network infrastructure that may be present in a number of homes. For example, the Home Phone Network Alliance (HPNA) allows users to network home computers by using the existing telephone and coaxial cable wiring within a home. HPNA-enabled devices utilize a different frequency spectrum than the spectrum used by faxes and phones. Instead of using existing telephone and coaxial wiring, the Homeplug® Power Alliance utilizes the existing electrical power wiring in a home to create a home network. In a Homeplug® network, all Homeplug®-enabled devices that are plugged into a wall outlet connected to a common electrical circuit may be wired together in a home network. One issue with Homeplug® is that the network bandwidth is susceptible to significant reduction due to large variations of the home electrical wiring and reactive loads in the outlets.\nAdditionally, problems arise in implementing network devices that correctly interact with other network devices. These problems may inhibit the deployment of newer devices that provide later-developed services in the presence of older (legacy) devices. The emerging Multimedia over Coax Alliance (MoCA) standard architecture impacts this problem in that (1) network behaviors dynamically assign a device the “Network Coordinator” (NC) role in order to optimize performance, (2) only the device in the NC role is known to be able to schedule traffic for all other nodes in the network, and (3) a full mesh network architecture may be formed between any device and its peers.\nWith many potential applications sharing the same digital network, various applications have to compete for a limited amount of bandwidth, which compounds the distribution problem. A bandwidth intensive application, such as a high-throughput download, may cause the degradation of other more important applications sharing the network. This outcome may be unacceptable when the other application requires a high quality of service.\nVarious solutions to solve this problem have been proposed, usually involving a high-level network controller or having high-level applications setting priority to data packets or data streams within the network. Moreover, intelligent network devices require high computational power, and are consequently more expensive than they need to be. Finally, complex network devices are impractical for home use, as most consumers do not have the sophistication or experience to configure a computer network."}
{"text": "The presently disclosed embodiments relate generally to layers that are useful in imaging apparatus members and components, for use in electrostatographic, including digital, apparatuses. More particularly, the embodiments pertain to an improved electrostatographic imaging member comprising an interfacial layer further comprising a semi-crystalline polyester resin to prevent light transmission to the substrate and thus significantly reduce “plywood effect,” a print quality defect.\nIn electrophotographic or electrostatographic printing, the charge retentive surface, typically known as a photoreceptor, is electrostatically charged, and then exposed to a light pattern of an original image to selectively discharge the surface in accordance therewith. The resulting pattern of charged and discharged areas on the photoreceptor form an electrostatic charge pattern, known as a latent image, conforming to the original image. The latent image is developed by contacting it with a finely divided electrostatically attractable powder known as toner. Toner is held on the image areas by the electrostatic charge on the photoreceptor surface. Thus, a toner image is produced in conformity with a light image of the original being reproduced or printed. The toner image may then be transferred to a substrate or support member (e.g., paper) directly or through the use of an intermediate transfer member, and the image affixed thereto to form a permanent record of the image to be reproduced or printed. Subsequent to development, excess toner left on the charge retentive surface is cleaned from the surface. The process is useful for light lens copying from an original or printing electronically generated or stored originals such as with a raster output scanner (ROS), where a charged surface may be imagewise discharged in a variety of ways.\nThe described electrostatographic copying process is well known and is commonly used for light lens copying of an original document. Analogous processes also exist in other electrostatographic printing applications such as, for example, digital laser printing or ionographic printing and reproduction where charge is deposited on a charge retentive surface in response to electronically generated or stored images.\nTo charge the surface of a photoreceptor, a contact type charging device has been used. The contact type charging device includes a conductive member which is supplied a voltage from a power source with a D.C. voltage superimposed with a A.C. voltage of no less than twice the level of the D.C. voltage. The charging device contacts the image bearing member (photoreceptor) surface, which is a member to be charged. The outer surface of the image bearing member is charged with the rubbing friction at the contact area. The contact type charging device charges the image bearing member to a predetermined potential. Typically the contact type charger is in the form of a roll charger such as that disclosed in U.S. Pat. No. 4,387,980, the relative portions thereof incorporated herein by reference.\nMultilayered photoreceptors or imaging members have at least two layers, and may include a substrate, a conductive layer, an optional undercoat layer (sometimes referred to as a “charge blocking layer” or “hole blocking layer”), an optional adhesive layer (sometimes referred to as an “interfacial layer”), a photogenerating layer (sometimes referred to as a “charge generation layer,” “charge generating layer,” or “charge generator layer”), a charge transport layer, and an optional overcoating layer in either a flexible belt form or a rigid drum configuration. In the multilayer configuration, the active layers of the photoreceptor are the charge generation layer (CGL) and the charge transport layer (CTL). Enhancement of charge transport across these layers provides better photoreceptor performance. Multilayered flexible photoreceptor members may include an anti-curl layer on the backside of the substrate, opposite to the side of the electrically active layers, to render the desired photoreceptor flatness.\nCoherent illumination is used in electrophotographic printing for image formation on photoreceptors. Unfortunately, the use of coherent illumination sources in conjunction with multilayered photoreceptors results in the “plywood effect,” also known as “interference fringe effect.” This defect consists of a series of dark and light interference patterns that occur when the coherent light is reflected from the interfaces that pervade multilayered photoreceptors. In organic photoreceptors, primarily the reflection from the undercoat layer or charge blocking layer/substrate interface (e.g., substrate surface) or the reflected light from the undercoat layer (or charge blocking layer)/charge generating layer interface account for the interference fringe effect. The effect can be eliminated if the strong undercoat layer surface reflection or the strong substrate surface reflection is eliminated or suppressed.\nWhile there have been attempts to reduce plywood effect in xerography, for example, as disclosed in U.S. Pat. Nos. 5,051,328, 5,089,908, 5,096,792, 5,139,907, 5,162,183, 5,210,548, 5,215,839, 5,302,485, 5,460,911, 5,635,324, 6,048,658, 6,214,514, 6,416,389, 6,582,872, and 7,314,812, hereby incorporated by reference, those methods were either too costly or had low efficiency in solving the interference problem.\nThus, there is a need for a low-cost, practically implementable and improved imaging layer that does not suffer from the above-described problems.\nConventional photoreceptors are disclosed in the following patents, a number of which describe the presence of light scattering particles in the undercoat layers: Yu, U.S. Pat. No. 5,660,961; Yu, U.S. Pat. No. 5,215,839; and Katayama et al., U.S. Pat. No. 5,958,638. The term “photoreceptor” or “photoconductor” is generally used interchangeably with the terms “imaging member.” The term “electrostatographic” includes “electrophotographic” and “xerographic.” The terms “charge transport molecule” are generally used interchangeably with the terms “hole transport molecule.”"}
{"text": "Transistors, in particular MOS transistors, like MOSFET or IGBT, can be used as electronic switches. In those applications a load path of the transistor is connected in series to the load, with the series circuit including the transistor and the load being connected between voltage supply terminals. The transistor can be switched on and off by applying a suitable control signal to its control input.\nIn many applications it is desired to measure the current through the load. For measuring the current through a load a sense transistor may be used. The sense transistor is coupled to the load transistor and may be operated in the same operation point as the load transistor. If the load transistor and the sense transistor are operated in the same operation point, a sense current flowing through the sense transistor is proportional to a current flowing through the load transistor. A proportionality factor between the load current and the sense current is dependent on the ratio between the active transistor areas (channel areas) of the sense transistor and the load transistor.\nThe proportionality factor is equal to the ratio between the active device areas, if the load transistor and the sense transistor have identical device characteristics. However, in reality the device characteristics of the load transistor and the sense transistor may differ. In one embodiment the threshold voltages of the two transistors may be different. This difference in the device characteristics may result from systematic differences in the processes that are used for manufacturing the load transistor and the sense transistor.\nUsually the load transistor and the sense transistor are integrated in a common semiconductor body (die, chip), where each of the load transistor and the sense transistor may include a number of transistor cells, that can be produced using identical processes. A first number of transistor cells is connected in parallel, thereby forming the load transistor, and a second number of transistor cells is connected in parallel, thereby forming the sense transistor. The transistor cells of the load and the sense transistor may have a common control terminal and one common first load terminal, but two different second load terminals. If the transistors are MOSFET the control terminal is a gate terminal, the first load terminal may be a drain terminal, and the second load terminals may be source terminals. Producing two different second load terminals (source terminals) may result in process variations that differently influence the threshold voltage of the load transistor cells and the sense transistor cells.\nFor these and other reasons there is a need for the present invention."}
{"text": "High-pressure high-temperature (HPHT) cells are used to form ultra-hard materials such as polycrystalline diamond (PCD) and polycrystalline cubic boron nitride (PCBN), which in turn are used in tools such as cutting tools and rock drilling tools. HPHT cells are used in HPHT presses such as, cubic presses, belt presses, and toroid presses. To form ultra-hard materials, HPHT presses often apply pressures in the range of 5 to 8 GPa and temperatures in the range of 1300 to 1650° C. For example, PCD may be sintered at 5 to 7 GPa and 1400 to 1500° C.\nThe formation of certain ultra-hard materials, such as thermally stable PCD, involves sintering at much higher temperatures. In particular, PCD formed using a carbonate catalyst may be sintered at a pressure greater than 6.5 GPa and a temperature greater than 2000° C. Additionally, binderless nano-polycrystalline PCD may be sintered at a pressure of about 15 GPa and a temperature of about 2300° C.\nSintering at these temperatures may be carried out using an HPHT cell including refractory materials that can withstand the high temperatures attained within the cell. For example, materials such as magnesium oxide (MgO) and sodium chloride (NaCl) have been used as pressure transfer media, and cubic zirconium oxide (ZrO2) having low thermal conductivity has been used as a thermal insulation sleeve in an HPHT cell. However, at pressures of about 8 GPa, NaCl begins to melt when the temperature exceeds 1600° C."}
{"text": "1. Field of the Invention\nThis invention relates to the field of integrated circuits, and in particular to the design of an interleaved transformer device with improved coupling characteristics.\n2. Description of Related Art\nA transformer comprises two coils of wire that are electromagnetically coupled. An AC signal applied to one of the coils of wire, nominally termed the primary coil, induces a corresponding AC signal on the other coil, nominally termed the secondary coil. The efficiency of the coupling between the coils determines the efficiency of the energy transfer from one coil to the other, as well as the overall noise sensitivity of the transformer.\nTransformers are generally used to provide electrical isolation between signals, often to allow for a change of voltage reference planes. Transformers are also often used to couple a signal to a resonant L-C circuit, one of the coils of the transformer providing the inductance, or a portion of the inductance. Radio-frequency (RF) transmitters typically include an RF transformer that couples an output stage to a resonant antenna stage.\nDiscrete transformers include two coils of wire that are wrapped around a core, which may be an air core, or a ferrite core, for improved coupling efficiency. The coils may be stacked, one atop the other on a common core, or may be wrapped one after the other on the common core, or may be formed by interleaving the wires and then wrapping the combination of the wires about the common core.\nStacked and interleaved transformers may be similarly constructed on multi-layer integrated circuits. Stacked coils can be implemented by forming conductive spirals of conductors on each of two layers, one atop the other. A transformer with interleaved coils is formed by laying the conductors of each coil adjacent to each other on the same conductive layer in a spiral pattern. A stacked core exhibits a high capacitance, because the two co-located spirals, one atop the other, effectively form two plates of a capacitor. This capacitive coupling between the coils introduces phase-shift and amplitude errors during the coupling process. In a transformer with interleaved coils, because the conductors are adjacent each other, rather than atop each other as in the stacked embodiment, the capacitive coupling is substantially less, thereby providing less phase-shift and less amplitude error.\nFIG. 1 illustrates an example layout of a coplanar interleaved transformer 100. Terminals 111 and 112 are the terminals of a first coil 110, and terminals 121 and 122 are the terminals of a second coil 120. Different cross-hatching is used for the conductors of each coil, for ease of illustration and understanding. Both coils are substantially on the same integrated circuit layer, preferably the heavy metallic layer, for minimal resistance losses. Interconnection segments 102 are located on a second circuit layer, to allow for insulated cross-overs of routing, as required. Connections between layers are illustrated by the circular areas 105, for ease of distinction in the illustrations. The coplanar interleaved transformer of FIG. 1 provides less capacitive coupling between coils than a stacked configuration, but consumes more area than a stacked configuration, and provides less efficient coupling than the stacked configuration.\nU.S. Pat. No. 5,543,773, xe2x80x9cTRANSFORMERS AND COUPLED INDUCTORS WITH OPTIMUM INTERLEAVING OF WINDINGSxe2x80x9d, issued Aug. 6, 1996 to Peter D. Evans and William J. B. Heffernan, teaches the interleaving of parallel segments of one of the coils, to achieve different turn-ratios, and is incorporated by reference herein. The referenced patent includes a coplanar transformer having a continuous primary coil spiral wound about a ferrite core, with interleaved segments of a secondary coil that are connected in parallel to effect a non-unity turns ratio. To improve the coupling between the primary and secondary coils, additional turns are conventionally used. For example, in the referenced patent, a 12:2 turns-ratio of actual turns is used to effect a 6:1 electrical turns-ratio; a 12:6 actual turns-ratio is used to effect a 2:1 turns-ratio, and so on. These additional actual turns increase the inductive coupling between the coils, while still maintaining the desired electrical turns-ratio that determines the ratio between input and output currents (n1*i1=n2*i2). Increasing the number of actual turns, however, increases the resistance, as well as the inductance, of each coil. In a coplanar transformer, increasing the number of actual turns substantially increases the area consumed by the transformer.\nIt is an object of this invention to provide a transformer structure that provides the relatively low capacitance of a conventional coplanar interleaved transformer, but with improved coupling efficiency and effectiveness. It is a further object of the invention to provide a coplanar transformer structure that provides improved coupling efficiency and effectiveness, without substantially increasing the inductance or resistance of the coils.\nThese objects and others are achieved by splitting each coil into two or more parallel segments, and interleaving the multiple segments of each coil with each other to form an interleaved coil that has a greater coupling efficiency than a non-segmented interleaved transformer. Because the multiple segments of each coil are connected in parallel, the resistance of the coil is decreased, or, the width of the conductor used to form the coil can be decreased while maintaining the same resistance. By reducing the width of the conductors used in a coplanar transformer, the additional parallel segments can be placed in substantially the same area as the conventional transformer, thereby maintaining the same inductance. To provide for maximum efficiency, each segment of each coil is embodied so as to have substantially equal length as each other segment of the coil."}
{"text": "1. Technical Field\nThe present invention relates to the use of gabapentin in a gastric retained dosage form. More specifically, the invention relates to the use of such dosage form to treat epilepsy and other disease states.\n2. Background\nGabapentin (1-(aminomethy)cyclohexaneacetic acid) is an anti-epileptic drug that is currently available in 100 mg, 300 mg and 400 mg hard shell capsule as well as 600 mg and 800 mg tablet dosage forms, with recommended dosing of 900 mg to 1800 mg total daily dose in three divided dosages. The oral bioavailability is dose-dependent, with approximately 60% bioavailability for a dose in the range of 300-400 mg, but with only 35% bioavailability for a dose of 1600 mg (Bourgeois, Epilepsia 36 (Suppl. 5):S1-S7 (1995); Gram, Epilepsia 37 (Suppl. 6):S12-S16 (1996)). The decrease in bioavailability with dose has been attributed to carrier-mediated absorption (Stewart, et al., Pharmaceutical Research 10(2):276-281 (1993).\nIn early work with rats, Vollmer, et al., Arzneim-Forsch/Drug Research 36(1, Nr. 5):781-892 (1986) found that the absorption site for gabapentin was the duodenum. The absorption of gabapentin occurs relatively slowly with the peak plasma concentration occurring approximately 2-6 hours after dosing (Bourgeois, supra). The elimination of gabapentin is exclusively through renal pathways (Chadwick; The Lancet 343:89-91 (1994); Vollmer, supra; Thomson, et al., Clin. Pharmacokinet. 23(3):216-230 (1992); and Riva, et al., Clin. Pharmacokinet. 31(6):470-493 (1996)) with reported half-lives of 5-7 hours (Chadwick, supra) and 6-7 hours (Gram, supra).\nA once- or twice-daily dosage form of gabapentin would be expected to improve compliance and therefore a controlled release dosage form has some distinct advantages over the conventional immediate release formulations. In addition, a controlled release dosage form would lower the maximum plasma concentration, and this may result in reduced side effects. Since gabapentin is absorbed high in the gastrointestinal tract, by means of a saturable transport mechanism, a gastric retained dosage form is particularly beneficial for delivery of gabapentin since the dosage form would be able to keep the drug in the region of absorption and show improved bioavailability by virtue of the slower release rate that avoids saturation of the carrier mediated transport of conventional dosages."}
{"text": "1. Field of the Instant Disclosure\nThe instant disclosure relates to a foldable keyboard; in particular, to a foldable keyboard comprising a folding support member having magnetic latch that can be used to retain the foldable keyboard in a folded position.\n2. Description of Related Art\nAs potable electronic devices, with touch panel, such as tablet computer and smartphone, becoming increasingly popular, user have become accustomed to providing input through typing on the touch screen, because most of compact electronic devices are not equipped with physical keyboards. However, while being intuitive and convenient for casual usage, the touching input method can only support limited typing speed. When a user needs to provide input of great volume, on screen input method through touch panes becomes impractical to. Thus, the compact foldable keyboards are provided to address this problem.\nTaiwan Patent NO. 1270807 disclosures a foldable keyboard that comprises two key portions: an upper cover and a lower cover. The upper cover and lower cover are configured to snap onto each other through a mechanical latch to provide a placing space, and the two key portions, while in the folded state, can place in the placing space. Moreover, the upper cover is configured to support an electronic device. However, the mechanical latch, which is made by plastic material, wears out relatively easily, thus reduce the durability of the product.\nTaiwan Patent NO. M419143 discloses a foldable keyboard comprising a left keyboard module and a right keyboard module assembled together through a connection module including a shaft, a subsidiary board, a pin axis and a plug. The right keyboard module and left keyboard module can rotate freely around the axis defined by the connection module. Thereby, the keyboard can be folded. However, the connection is complicated with a amount of components, which translates to higher production cost and manufacture complexity.\nTo address the above issues, the inventions strive via industrial experience and academic research to present the instant disclosure, which can effectively improve the limitations described above."}
{"text": "1. Field of the Invention\nThe present application relates to an over-current protection device, and more particularly to an over-current protection device having positive temperature coefficient (PTC) characteristic. The novel over-current protection device is configured to replace a traditional circuit breaker.\n2. Description of Related Art Including Information Disclosed Under 37 CFR 197 and 137 CFR 1.98\nIn recent years, super-thin portable electronic apparatuses have been increasingly developed, resulting in the demand of lightweight, thin and large capacity batteries. To meet the demand of high voltage and large capacity, a battery may use a large capacity cell or a plurality of cells in parallel and/or series connection. In over-current and over-heat protection to batteries, circuit breakers in response to temperature are usually employed nowadays.\nAn over-current breaker usually uses bi-metal plates, and the bi-metal plates constitute a normally closed micro-switch. When the breaker has current flowing therein, the bi-metal plates are heated and bent and the bi-metal plates bent more intensively as current increases. If over-current occurs, extremely large current induces over-bending of the bi-metal plates, the micro-switch changes to and sustain in an open state to sever the current. As a result, the breaker prevents the circuit from damage which may be caused by over-current. When the over-current or over-temperature event has gone, the bi-metal plates cool down and return to be of original shapes to rebuild a conductive path.\nFIG. 1 shows an appearance of a bi-metal breaker 5. The first metal plate 6 and the second metal plate 7 extend outwardly from the insulating body 8. The first metal plate 6 and the second metal plate 7 in the insulating layer 8 form a normally closed micro-switch, whereas the first metal plate 6 and the second metal plate 7 out of the insulating layer 8 serve as soldering, or welding interfaces.\nAlthough the bi-metal design has been widely used in breakers, the bi-metal plates have to be made by precision machinery and need precision manufacturing techniques. As a result, the bi-metal breakers are usually costly. Therefore, it is highly demanded to generate a breaker with high reliability and stability for safety, and consider how to simplify manufacturing process and obtain low cost."}
{"text": "In recent years, nucleic acid hybridization has become an increasingly important means of identifying, measuring and detecting the presence of particular nucleic acids in a given sample. Thus, for example, the fields of medical diagnostics, environmental and food testing, and forensics have all benefitted from the use of nucleic acid hybridization as a rapid, simple and extraordinarily accurate way of testing for the presence or absence of given biological contaminants or microorganisms in a sample.\nMost nucleic acid hybridization schemes have features in common. One such typical feature is the use of single-stranded nucleic acid probes (or denatured double-stranded probes) having a defined or known nucleotide sequence. Probe molecules may be derived from biological sources, such as genomic DNA or RNA, or may be enzymatically synthesized, either in a prokaryotic or eukaryotic host cell or in vitro. Presently, most nucleic acid probes in common use are oligonucleotide probes made using chemical synthetic methods (“synthetic oligonucleotides”). One such synthetic method is automated sequential addition of 3′-activated, protected nucleotides to the 5′ end of a growing, solid phase-bound oligonucleotide chain, followed by cleavage of the completed oligonucleotide from the support and deprotection. See, e.g., Eckstein, Oligonucleotides & Analogues: A Practical Approach (1991).\nSynthetic oligonucleotides for use as hybridization probes are typically deoxyribonucleotides having a nucleotide sequence complementary to a nucleotide sequence of the nucleic acid to be detected. DNA oligonucleotides are classically preferred for a number of reasons. Among these is the greater stability DNA has to enzymatic hydrolysis upon exposure to common samples, due to the almost ubiquitous presence in samples of various RNAses. RNA is also known to be less chemically stable than DNA, e.g., RNA degradation is facilitated by the presence of base, heavy metals. And compared to RNA, DNA is less prone to assume stable secondary structures under assay conditions. Such secondary structures can render oligonucleotides unavailable for inter-molecular hybridization. Nevertheless, RNA oligonucleotides may be used, even though they are less preferred.\nNucleic acid hybridization exploits the ability of single-stranded nucleic acids to form stable hybrids with corresponding regions of nucleic acid strands having complementary nucleotide sequences. Such hybrids usually consist of double-stranded duplexes, although triple-stranded structures are also known. Generally speaking, single-strands of DNA or RNA are formed from nucleotides containing the bases adenine (A), cytosine (C), thymidine (T), guanine (G), uracil (U), or inosine (I). The single-stranded chains may hybridize to form a double-stranded structure held together by hydrogen bonds between pairs of complementary bases. Generally, A is hydrogen bonded to T or U, while G or I is hydrogen bonded to C. Along the double-stranded chain, classical base pairs of AT or AU, TA or UA, GC, or CG are present. Additionally, some mismatched base pairs (e.g., AG, GU) may be present. Under appropriate hybridization conditions, DNA/DNA, RNA/DNA, or RNA/RNA hybrids can form.\nBy “complementary” is meant that the nucleotide sequences of corresponding regions of two single-stranded nucleic acids, or two different regions of the same single-stranded nucleic acid, have a nucleotide base composition that allows the single strands to hybridize together in a stable double-stranded hydrogen-bonded region under stringent hybridization conditions. When a contiguous sequence of nucleotides of one single stranded region is able to form a series of “canonical” hydrogen-bonded base pairs with an analogous sequence of nucleotides of the other single-stranded region, such that A is paired with U or T, and C is paired with G, the nucleotides sequences are “perfectly” complementary.\nThe extreme specificity of nucleic acid hybridization, which under some circumstances can allow the discrimination of nucleic acids differing by as little as one base, has allowed the development of hybridization-based assays of samples containing specific microorganisms, nucleic acids bearing given genetic markers, tissue, biological fluids and the like. Such assays are often able to identify nucleic acids belonging to particular species of microorganisms in a sample containing other, closely-related species. Nucleic acid hybridization assays can also specifically detect or identify certain individuals, or groups of individuals, within a species, such as in the forensic use of RFLP (restriction fragment length polymorphism) and PCR (polymerase chain reaction) testing of samples of human origin.\nThe use of oligonucleotides as a diagnostic tool in nucleic acid hybridization testing often involves, but need not involve, the use of a reporting group or “label” which is joined to the oligonucleotide probe molecule, or both the probe and the target. Such a reporter group moiety may include, for example, a radioisotope, chemiluminescent or fluorescent agent, or enzyme joined to the oligonucleotide. The label is employed to render the probe capable of detection, particularly when the probe is hybridized to the target nucleic acid.\nThe majority of assay methods employing nucleic acids utilize a physical separation step in order to separate the probe:analyte hybrid from unbound probe. These assay methods are called “heterogeneous” assays. In nucleic acid hybridization assays, an analyte molecule is the target nucleic acid species sought to be detected, quantitated and/or identified. A “hybrid” is a partly or wholly double-stranded nucleic acid comprising two single-stranded nucleic acids, such as a probe and a target nucleic acid, having a region of complementarity resulting in intermolecular hydrogen bonding under assay and/or amplification conditions.\nAssay methods utilizing a physical separation step include methods employing a solid-phase matrix, such as glass, minerals or polymeric materials, in the separation process. The separation may involve preferentially binding the probe:analyte complex to the solid phase matrix, while allowing the unassociated probe molecules to remain in a liquid phase. Such binding may be non-specific, as, for example, in the case of hydroxyapatite, or specific, for example, through sequence-specific interaction of the target nucleic acid with a “capture” probe which is directly or indirectly immobilized on the solid support. In any such case, the amount of probe remaining bound to the solid phase support after a washing step is proportional to the amount of analyte in the sample.\nAlternatively, the assay may involve preferentially binding the unhybridized probe while leaving the hybrid to remain in the liquid phase. In this case the amount of probe in the liquid phase after a washing step is proportional to the amount of analyte in the original sample. When the probe is a nucleic acid or oligonucleotide, the solid support can include, without limitation, an adsorbent such as hydroxyapatite, a polycationic moiety, a hydrophobic or “reverse phase” material, an ion-exchange matrix, such as DEAE, a gel filtration matrix, or a combination of one or more of these solid phase materials. The solid support may contain one or more oligonucleotides, or other specific binding moiety, to capture, directly or indirectly, probe, target, or both. In the case of media, such as gel filtration, polyacrylamide gel or agarose gel, the separation is not due to binding of the oligonucleotide but is caused by molecular sieving of differently sized or shaped molecules. In the latter two cases, separation may be driven electrophoretically by application of an electrical current through the gel causing the differential migration through the gel of nucleic acids of different sizes or shapes, such as double-stranded and single-stranded nucleic acids.\nA heterogeneous assay method may also involve binding the probe to a solid-phase matrix prior to addition of a sample suspected of containing the analyte of interest. The sample can be contacted with the label under conditions which would cause the desired nucleic acid to be labeled, if present in the sample mixture. The solid phase matrix may be derivatized or activated so that a covalent bond is formed between the probe and the matrix. Alternatively, the probe may be bound to the matrix through strong non-covalent interactions, including, without limitation, the following interactions: ionic, hydrophobic, reverse-phase, immunobinding, chelating, and enzyme-substrate. After the matrix-bound probe is exposed to the labeled nucleic acid under conditions allowing the formation of a hybrid, the separation step is accomplished by washing the solid-phase matrix free of any unbound, labeled analyte. Conversely, the analyte can be bound to the solid phase matrix and contacted with labeled probe, with the excess free probe washed from the matrix before detection of the label.\nYet another type of assay system is termed “homogeneous assay.” Homogenous assays can generally take place in solution, without a solid phase separation step, and commonly exploit chemical differences between the free probe and the analyte:probe complex. An example of an assay system which can be used in a homogenous or heterogeneous format is the hybridization protection assay (HPA) disclosed in Arnold, et al., U.S. Pat. No. 5,283,174, in which a probe is linked to a chemiluminescent moiety, contacted with an analyte and then subjected to selective chemical degradation or a detectable change in stability under conditions which alter the chemiluminescent reagent bound or joined to unhybridized probe, without altering the chemiluminescent reagent bound or joined to an analyte:probe conjugate. Subsequent initiation of a chemiluminescent reaction causes the hybrid-associated label to emit light. This patent enjoys common ownership with the present application and is expressly incorporated by reference herein.\nCompetition assays, in which a labeled probe or analyte competes for binding with its unlabeled analog, are also commonly used in a heterogeneous format. Depending on how the system is designed, either the amount of bound, labeled probe or the amount of unbound, labeled probe can be correlated with the amount of analyte in a sample. However, such an assay can also be used in a homogeneous format without a physical separation step, or in a format incorporating elements of both a homogeneous and a heterogeneous assay.\nThe assay methods described herein are merely illustrative and should not be understood as exhausting the assay formats employing nucleic acids known to those of skill in the art.\nNucleic acid hybridization has been utilized in methods aimed at using oligonucleotides as therapeutic agents to modify or inhibit gene expression within living organisms. In an example of such utilization, oligonucleotide “antisense” agents can be targeted specifically to an mRNA species encoding a deleterious gene product, such as a viral protein or an oncogene. See, e.g.; Zamecnik and Stephenson, 75 Proc. Nat'l Acad. Sci. (USA), 280-284 (1978); Stephenson and Zamecnik, 75 Proc. Nat'l Acad. Sci. (USA), 285-288 (1978); and Tullis, U.S. Pat. No. 5,023,243. Although Applicant does not wish to be bound by theory, it is thought that the RNA:DNA duplex which results from the binding of the antisense oligonucleotide to RNA targets may serve as a substrate for RNAse H, an RNA-degrading enzyme present in most cells and specific for RNA contained in an RNA:DNA duplex. According to this model, the target RNA molecule is destroyed through hybridization to the antisense oligonucleotide. Variations of this general strategy exist, wherein, for example, the oligonucleotide has a structure conferring an enzymatic activity on the oligonucleotide, such as the RNAse activity of so-called ribozymes. See, e.g., Goodchild, PCT Publication No. WO93/15194.\nBecause therapeutic antisense oligonucleotides are primarily designed to function in vivo formulations for the delivery of such agents must not significantly inhibit normal cellular function. Thus, nuclease inhibitors, which can sometimes be included in in vitro diagnostic tests to prevent oligonucleotide degradation, are not suitable for use in vivo. This fact has resulted in the design of various oligonucleotides modified at the internucleotide linkage, at the base or sugar moieties, or at combinations of these sites to have greater nuclease resistance than unmodified DNA.\nThus, a number of oligonucleotide derivatives have been made having modifications at the nitrogenous base, including replacement of the amino group at the 6 position of adenosine by hydrogen to yield purine; substitution of the 6-keto oxygen of guanosine with hydrogen to yield 2-amino purine, or with sulphur to yield 6-thioguanosine, and replacement of the 4-keto oxygen of thymidine with either sulphur or hydrogen to yield, respectively, 4-thiothymidine or 4-hydrothymidine. All these nucleotide analogues can be used as reactants for the synthesis of oligonucleotides. See, e.g., Oligonucleotides and Analogues: A Practical Approach, supra. Other substituted bases are known in the art. See, e.g., Cook, et al., PCT Publication No. WO92/02258, entitled Nuclease Resistant, Pyrimidine Modified Oligonucleotides that Detect and Modulate Gene Expression, which is incorporated by reference herein. Base-modified nucleotide derivatives can be commercially obtained for oligonucleotide synthesis.\nSimilarly, a number of nucleotide derivatives have been reported having modifications of the ribofuranosyl or deoxyribofuranosyl moiety. See, e.g., Cook et al., PCT Publication No. WO94/19023 entitled Cyclobutyl Antisense Oligonucleotides, Methods of Making and Use Thereof; McGee, et al., PCT Publication No. WO94/02501 entitled Novel 2′-O-Alkyl Nucleosides and Phosphoramidites Processes for the Preparation and Uses Thereof; and Cook, PCT Publication No. WO93/13121 entitled Gapped 2′-modified Oligonucleotides. These three publications are incorporated by reference herein.\nMost oligonucleotides comprising such modified bases have been formulated with increased cellular uptake, nuclease resistance, and/or increased substrate binding in mind. In other words, such oligonucleotides are described as therapeutic gene-modulating agents.\nNucleic acids having modified nucleotide residues exist in nature. Thus, depending on the type or source, modified bases in RNA can include methylated or dimethylated bases, deaminated bases, carboxylated bases, thiolated bases and bases having various combinations of these modifications. Additionally, 2′-O-alkylated bases are known to be present in naturally occurring nucleic acids. See, Adams, The Biochemistry of the Nucleic Acids, 7,8 (11th ed. 1993)."}
{"text": "A failover cluster is a group of server computers that work together to maintain high availability of applications and services. The server computers may also be referred to as “servers” or as “nodes”. If one of the nodes fails, then another node in the cluster can take over its workload with little or no downtime, a process which is known as “failover”.\nIn some computing architectures, a group of computers that share a network resource database and have a common security policy are called a “domain”. A domain may have a domain controller, sometimes called a “primary domain controller” (PDC) which is a computer that manages the resource and user access for the entire domain. One or more backup domain controllers may also be present. Other computers in the domain are workstations or servers that provide computational resources to domain users, after domain authentication or another authentication protocol establishes that the users have permission to access the computers in the domain. A domain administrator may manage access using administrative software and security policies."}
{"text": "Electric power distribution panelboards and switchboards are currently available which accept circuit breakers and electric switches by means of a plug-on connection to the bus bars arranged therein. U.S. Pat. No. 4,744,003 describes one such panelboard arrangement. When electric switches or circuit breakers are to be installed on-site within such an operating panelboard, the electric power to the panelboard must be temporarily discontinued and the panelboard disassembled before such electric switch or circuit breaker can be installed. Where modification must be made to the panelboard interior in order to accept the electric switch or circuit breaker, some time is required before the power can be turned ON. This disruption in electric power within an industrial environment could cause scheduling problems both with respect to the related equipment as well as to operating personnel. To provide adequate electrical isolation between the different phases of a multi-phase panelboard or switchboard assembly, the bus conductors that transfer operating power from the electrical distribution circuits to the various switches connected with the panelboards and switchboards, large bus conductors are supported within the panelboard and switchboard interiors and heavy plates of insulative material are used to both support and electrically isolate the bus conductors. When the panelboards and switchboards are operated at near capacity levels, the heavy insulative plates thermally insulate the bus conductors and interfere with the convection air flow. Since the panelboard and switchboard enclosures are required to be closed in accordance with electrical industries standards, the temperature within the enclosures can rise substantially above ambient resulting in eventual damage to the panelboard components.\nIn accordance with the state-of-the-art panelboard and switchboard enclosures, a wide variety of electric switch designs are inventoried in order to meet the ampere requirements of varied electrical distribution circuits. The increased ampere-rated electric switches require larger components in order to accommodate the higher ampacity and to prevent the electric switches from overheating during long term steady state operating conditions.\nIt would be economically advantageous, to provide a single electric switch design that can operate over a wide range of ampere ratings and thereby eliminate the large inventory requirement of several electric switches of various component sizes. It would be further advantageous, to control the thermal environment within the panelboard and switchboard enclosure to safely and economically operate the switches at reduced operating temperatures.\nU.S. Pat. No. 4,752,233, which Patent is incorporated herein for reference purposes, describes an adapter module for interfacing between an electric switch or circuit breaker and a panelboard or switchboard housing. The adapter module contains two compartments which allow two such circuit breakers, electric switches or a combination thereof to be plugged onto the panelboard or switchboard bus bars at one time. The adapter module also allows one electric switch or circuit breaker to be connected within one compartment while reserving the other compartment for add-on capability after the adapter is installed within the panelboard-switchboard enclosure.\nU.S. Pat. Nos. 4,675,481; 4,785,145; 4,778,959; 4,778,961 each describe various electric switch configurations that can be received within the adapter module during factory or on-site installation. Since the components of each of the above described electric switches can be automatically assembled to reduce assembly costs, it would be further economically advantageous to assemble such switches in a panelboard or switchboard enclosure without incurring additional labor and component costs.\nFurther, when an electric switch or circuit breaker is required to be added to existing panelboard and switchboard enclosures, it would be economically advantageous to obtain such electric switch or circuit breaker in the form of a kit that can be quickly assembled within the panelboard and switchboard enclosure without having to completely disassemble the associated panelboard and switchboard housing.\nOne purpose of this invention accordingly is to provide a panelboard-switchboard housing that easily and quickly accepts multiple electric switch-circuit breaker enclosures without requiring long-time equipment shut-down. A further purpose of this invention is to provide an electric switch or circuit breaker modular enclosure that allow additional electric switches and circuit breakers to be field-installed at a minimum expense of time and labor.\nAn additional purpose of this invention is to provide a panelboard-switchboard housing that is designed to control the temperature rise of the bus conductors, electric switches and circuit breakers mounted within the housing."}
{"text": "The present invention relates to a tongue of an inflatable belt device for protecting an occupant of a vehicle such as an automobile in the event of a vehicle collision. More particularly, the present invention relates to a tongue to which a shoulder belt and a lap belt are both connected.\nAs an example of this type of inflatable belt device, there is a device disclosed in Japanese Patent Unexamined Publication No. H5-85301 (incorporated by reference herein). FIG. 13 is a perspective view of an inflatable belt device disclosed in the publication.\nAs shown in FIG. 13, the inflatable belt device 120 comprises an inflatable belt 122 configured as a shoulder belt extending diagonally from the right side to the left side of a vehicle occupant. The device also includes a webbing 122a connected to the inflatable belt 122, a lap belt 123 extending from the right side to the left side of the occupant, and a buckle 124 fixed to, for example, a vehicle floor. A tongue 125 to be inserted into and latched with the buckle 124 when the occupant wears the inflatable belt is also provided. An intermediate guide 126 is provided for guiding the webbing 122a. \nThe webbing 122a is a normal webbing which is the same as a typical conventional seat belt and is slidably hung by the intermediate guide 126. The end of the webbing 122a is connected to a seat belt retractor (ELR) 127, which has an emergency locking mechanism and is fixed to the vehicle body.\nOne end of the inflatable belt 122 is connected to the webbing 122a by sewing or the like and the other end of the inflatable belt 122 is connected to the tongue 125.\nThe lap belt 123 is a normal webbing which is the same as a typical conventional seat belt and of which one end is connected to the tongue 125 and the other end is connected to a seat belt retractor (ELR) 128, which is fixed to the vehicle body. A gas generator 129, which is actuated to generate high pressure gas in emergency situations, e.g., a collision, is connected to the buckle 124.\nThe tongue 125 and the buckle 124 have passages allowing gas from the gas generator 129 to flow into the inflatable belt 122.\nThe inflatable belt 122 comprises a bag formed in a long envelope shape having a large width and a cover. The bag is folded into a band-like shape having a small width and is covered by the cover. In the normal state, the bag is held in the band-like shape. The inflatable belt device 120 can be used in the same manner as the normal seat belt device. As the gas generator 129 is actuated in the event of a vehicle collision, the inflatable belt 122 is inflated, thereby protecting the vehicle occupant.\nAn example of structure for connecting the lap belt to the tongue is disclosed in Japanese Patent Unexamined Publication No. H11-334513 (incorporated by reference herein). In particular, FIG. 5 of the publication shows the structure. This structure is configured so that a lap belt anchor is attached to a tongue plate and a lap belt is connected to the lap belt anchor. The lap belt anchor is provided with an opening through which the lap belt is inserted. After an end portion of the lap belt is inserted into the opening, the lap belt is folded double along an edge of the opening so that the end portion is superposed on a portion of the lap belt following the end portion. The superposed portions of the lap belt are sewn together by stitching.\nIn one example of structure of installing an inflatable belt device comprising a tongue, a shoulder belt, a lap belt, and a retractor to an automobile, the retractor is arranged inside a B-pillar and is covered by a pillar cover. The shoulder belt and the lap belt are passed through belt through holes of the pillar cover.\nIn the inflatable belt device in which the lap belt anchor attached to the tongue plate is connected to the lap belt by sewing just like the device disclosed in Japanese Patent Unexamined Publication No. H11-334513, the lap belt anchor is attached to the tongue plate and the lap belt is connected to the lap belt anchor in the manufacturing process of the inflatable belt device. In an inflatable belt device of this type and in the aforementioned installation structure in which the retractor is arranged inside a B-pillar, the shoulder belt and the lap belt are required to be previously passed through the belt through holes of the pillar cover before the assembly line of the automobile, because it is impossible to sew the lap belt by using a sewing machine during the assembly line of the automobile.\nThe inflatable belt device assembled to the pillar cover is bulky, thus increasing the cost for carrying it to the assembly line of the automobile.\nTherefore, it is an object of the present invention to provide a tongue of an inflatable belt device which allows a lap belt anchor to be easily attached to a body of the tongue, whereby the inflatable belt device can be carried solely (that is, without being assembled with a vehicle interior cover such as a pillar cover) to the assembly line of an automobile.\nAccording to one embodiment of the present invention, a tongue of an inflatable belt device is provided. The tongue includes a tongue body which can be latched with a buckle and to which a shoulder belt including an inflatable belt is connected. The tongue also includes a lap belt anchor to be coupled to the tongue body, and at least one insertion member of the tongue body for engaging the lap belt anchor. A stopper is disposed on the tongue body for preventing the lap belt anchor engaged with the insertion member from coming off the insertion member. The stopper can be elastically displaced or deformed from a prevention position where the stopper prevents the lap belt anchor from coming off the insertion member to an escape position where the stopper allows the lap belt anchor to be engaged with the insertion member.\nAccording to the present invention, the tongue of the inflatable belt, the lap belt anchor may be coupled to the tongue body by engaging the lap belt anchor with the insertion member of the tongue body. As a result, the lap belt anchor and the tongue body may be separately carried to the assembly line of automobiles. After passing the shoulder belt and the lap belt with the lap belt anchor into openings of a pillar cover, the lap belt anchor can be coupled to the tongue body during assembly of the automobile. Therefore, the inflatable belt device can be carried without being assembled with a pillar cover to the assembly line for production of an automobile.\nWhen the lap belt anchor is engaged with the insertion member, the stopper is in the escape position so that the lap belt anchor can be smoothly engaged with the insertion member. As the lap belt anchor is fully engaged with the insertion member, the stopper is elastically returned to its prevention position, thereby securely coupling the lap belt to the tongue body.\nAccording to another embodiment of the present invention, a tongue of an inflatable belt device includes a tongue body which can be latched with a buckle and to which a shoulder belt including an inflatable belt is connect. A lap belt anchor may be coupled to the tongue body. The tongue body is provided with at least one insertion opening for the lap belt anchor and the lap belt anchor is provided with at least one hook to be inserted into the insertion opening in order to engage with the tongue body. The hook may be elastically displaced or deformed in order to pass through the insertion opening. The hook is returned to its original position or its original configuration after passing through the insertion opening in order to engage with the tongue body.\nThe lap belt anchor may be coupled to the tongue body by inserting the lap belt anchor having the insertion opening. The hook may be displaced or deformed to pass through the insertion opening when the lap belt anchor is inserted into the insertion opening, so that the lap belt anchor can be inserted deeply into the insertion opening. After the hook passes through the insertion opening, the hook is returned to its holding position or holding configuration, thereby securely coupling the lap belt anchor to the tongue body.\nAccording to another embodiment of the present invention, a tongue of an inflatable belt device may include a tongue body which can be latched with a buckle and to which a shoulder belt composed of an inflatable belt may be connected. A lap belt through hole is formed in the tongue body. A lap belt is provided, a front end portion of the lap belt passes through the lap belt through hole and is held by an edge portion of the lap belt through hole. The lap belt is provided at its front end with a hoop portion. The front end of the lap belt is held by an edge portion of the lap belt through hole with a bar member inserted into the hoop portion.\nAccording to another embodiment of the present invention, the tongue of the inflatable belt device may be configured so that after the front end portion of the lap belt passes through the through hole of the tongue body, the bar member is inserted into the hoop portion of the front end portion of the lap belt in order to hold the front end portion of the lap belt at the edge portion of the through hole. Both ends of the bar member may be engaged by edge portions of the through hole. Alternatively, the hoop portion may be engaged by the edge portions of the through hole by increasing the thickness of the hoop portion due to the insertion of the bar member into the hoop portion.\nAccording to yet another embodiment of the present invention, bar member can be elastically deformed such that the bar member in the state inserted into the hoop portion can pass through said lap belt through hole, and can be elastically returned to be held by the edge portion of the lap belt through hole after passing through the lap belt through hole.\nAccording to the present invention, the bar member may be inserted into the hoop portion of the front end portion of the lap belt. The hoop portion is then passed through the lap belt through hole. The bar member elastically deforms to allow the hoop portion to pass through the lap belt through hole. After the hoop portion is passed through the lap belt through hole, the bar member elastically returns to its original form and engages the edge portion of the lap belt through hole.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory only, and are not restrictive of the invention as claimed."}
{"text": "It is not possible to change the arrangement of devices provided in a mobile terminal apparatus in accordance with the direction in which it is held.\nTherefore, there is a risk that the devices may not be used in an appropriate position."}
{"text": "This invention relates to a process for the production of hydrofluoroalkanes, particularly 1,1,1,2-tetrafluoroethane and pentafluoroethane.\nSeveral processes have been proposed for the production of 1,1,1,2-tetrafluoroethane, otherwise known as HFC 134a, and pentafluoroethane, otherwise known as HFC 125 which are employed as or as components of replacements for chiorofluorocarbons in the many applications in which chlorofluorocarbons are employed. Amongst such processes is the fluorination of the corresponding chlorine-containing starting material by reacting the starting material with hydrogen fluoride in the liquid or the vapour phase, usually in the presence of a fluorination catalyst.\nThus it has been proposed in United Kingdom Patent Specification No. 1,589,924 to produce HFC 134a by the vapour phase fluorination of 1,1,1-trifluoro-2-chloroethane (HCFC 133a) which is itself obtainable by the fluorination of trichioroethylene as described in United Kingdom Patent Specification No. 1,307,224.\nThe formation of HFC 134a as a minor product of the fluorination or trichloroethylene is described in United Kingdom Patent Specification No 819,849, the major reaction product being HCFC 133a.\nMore recently, processes for the production of HFC 134a from trichloroethylene based on a combination of the reaction of trichloroethylene with hydrogen fluoride to produce HCFC 133a and the reaction of HCFC 133a with hydrogen fluoride to produce HFC 134a have been proposed.\nIn WO 90/08755, the contents of which are incorporated herein by reference, there is described the conversion of trichloroethylene to HFC 134a wherein the two reactions steps are carried out in a single reaction zone with recycle of part of the product stream containing HCFC 133a.\nIn EP 0 449 614, the contents of which are also incorporated herein by reference, there is described a process for the manufacture of HFC 134a which comprises the steps of:\n(A) contacting a mixture of trichloroethylene and hydrogen fluoride with a fluorination catalyst under superatmospheric pressure at a temperature in the range from about 200xc2x0 C. to about 400xc2x0 C. in a first reaction zone to form a product containing 1,1,1-trifluoro-2-chloroethane and hydrogen chloride together with unreacted starting materials,\n(B) passing product of step A together with hydrogen fluoride to a second reaction zone containing a fluorination catalyst at a temperature in the range from about 280xc2x0 C. to about 450xc2x0 C. but higher than the temperature in step A to form a product containing 1,1,1-trifluoro-2-chloroethane, 1,1,1,2-tetrafluoroethane and hydrogen chloride,\n(C) treating product of step B to separate 1,1,1,2-tetrafluoroethane and hydrogen chloride from 1,1,1-trifluoro-2-chloroethane and unreacted hydrogen fluoride, and\n(D) feeding 1,1,1-trifluoro-2-chloroethane obtained from step C together with trichloroethylene and hydrogen fluoride to said first reaction zone (step A).\nIn EP 0 449 617, the contents of which are also incorporated herein by reference, there is described a process for the production of BFC 134a which comprises the steps of:\n(A) contacting a mixture of 1,1,1-trifluoro-2-chloroethane and hydrogen fluoride with a fluorination catalyst at a temperature in the range from about 280xc2x0 C. to about 450xc2x0 C. in a first reaction zone to form a product containing 1,1,1,2-tetrafluoroethane and hydrogen chloride together with unreacted starting materials,\n(B) passing product of step A together with trichloroethylene to a second reaction zone containing a fluorination catalyst at a temperature in the range from about 200xc2x0 C. to about 400xc2x0 C. but lower than the temperature in step A to form a product containing 1,1,1-trifluoro-2-chloroethane, 1,1,1,2-tetrafluoroethane, hydrogen chloride and unreacted trichloroethylene and hydrogen fluoride,\n(C) treating product of step B to separate 1,1,1,2-tetrafluoroethane and hydrogen chloride from 1,1,1-trifluoro-2-chloroethane, unreacted trichioroethylene and hydrogen fluoride, and\n(D) feeding 1,1,1-trifluoro-2-chloroethane obtained from step C together with hydrogen fluoride to said first reaction zone (step A).\nHowever, a problem which is encountered with processes for the production of 1,1,1,2-tetrafluoroethane based on the hydrofluorination of 1-chloro-2,2,2-trifluoroethane and/or trichioroethylene, is that the conversion of 1-chloro-2,2,2-trifluoroethane to 1,1,1,2-tetrafluoroethane is equilibrium limited, there being a maximum conversion of 1-chloro-2,2,2-trifluoroethane to 1,1,1,2-tetrafluoroethane of only about 20% under typical operating conditions.\nIt has also been proposed to produce pentafluoroethane (HFC 125) by the catalysed fluorination with hydrogen fluoride of chlorotetrafluoroethane (HCFC 124) and/or dichlorotrifluoroethane (HCFC 123) which are themselves obtainable by the fluorination of perchioroethylene with hydrogen fluoride.\nThe present invention resides in a process for the production of hydrofluoroalkanes, particularly 1,1,1,2-tetrafluoroethane and pentafluoroethane from hitherto unused starting materials, which process in the case of production of 1,1,1,2-tetrafluoroethane is not subject to the aforementioned equilibrium limitation problem.\nAccording to the present invention there is provided a process for the production of a hydrofluoroalkane which comprises contacting a hydrochlorofluoroethane having the formula CClXYCFHZ or a(hydro)chlorofluoroethene having the formula CClAxe2x95x90CFZ in which X and Y are each independently chlorine or fluorine, Z is chlorine or hydrogen and A is chlorine or fluorine provided that where each of X and Y is fluorine then Z is hydrogen in the vapour phase with hydrogen fluoride and a fluorination catalyst and recovering a hydrofluoroalkane from the resulting products.\nIn a particular embodiment of the process for producing 1,1,1,2-tetrafluoroethane, the hydrochlorofluoroethane has the formula CClXYCFH2 and the (hydro)chlorofluoroethene has the formula CClAxe2x95x90CFH.\nWe have found that the product gases from the process for producing 1,1,1,2-tetrafluoroethane comprise a greater molar proportion of 1,1,1,2-tetrafluoroethane than is obtained when 1-chloro-2,2,2-trifluoroethane is used as the starting material.\nThe starting materials for the process are CCl3CFH2, CCl2FCFH2, CClF2CFH2, CCl2FCClFH, CCl3CHFCl, CCl2xe2x95x90CFH, CClFxe2x95x90CFH, CCl2xe2x95x90CFCl and CClFxe2x95x90CFCl. We prefer to employ CCl2xe2x95x90CFH or CCl2FCFH2, especially CCl2xe2x95x90CFH for the production of 1,1,1,2-tetrafluoroethane and CCl2FCClFH or CClFxe2x95x90CFCl for the production of pentafluoroethane since these materials are more readily available.\nProcesses for the production of the starting materials of the present invention are known. Thus for example CCl2xe2x95x90CFH may be produced from trichloroethylene, as described in the Journal of Organic Chemistry 28, 112 (1963), or from tetrachloroethane as described in EP 537560.\nSuitable fluorination catalysts are those which yield the desired hydrofluoroalkane as a product of the reaction with a yield of greater than 20%, preferably greater than 25%, based on the starting material processed and include catalysts based on chromia or chromium oxyfluoride, and the fluorides or oxyfluorides of other metals, for example magnesium and aluminium. Activity promoting amounts of other metals, for example zinc and nickel may also be present; we particularly prefer to employ a catalyst comprising zinc on chromia as described fully in published European Patent Application No. 502605, the contents of which are incorporated herein by reference.\nThe relative proportion of hydrogen fluoride to starting material which is employed may vary within wide limits although it is generally preferred to employ a stoichiometric excess of hydrogen fluoride. The stoichiometrically required molar ratio depends upon the particular starting material. Where the starting material is the preferred 1,1-dichloro-2-fluoroethene, the stoichiometrically required molar ratio of hydrogen fluoride to 1,1-dichloro-2-fluoroethene is 3:1. The molar ratio of hydrogen fluoride to the starting material, for example 1,1-dichloro-2-fluoroethene, will usually be at least 2:1, and preferably is at least 4:1 and especially at least 6:1 and substantially greater excesses of hydrogen fluoride, for example up to 50:1, may be employed if desired.\nThe temperature at which the process is carried out is preferably at least 180xc2x0 C. and more preferably at least 200xc2x0 C. or 220xc2x0 C. but may be significantly lower than the temperatures typically employed for the conversion of 1-chloro-2,2,2-trifluoroethane to 1,1,1,2-tetrafluoroethane. Preferably the temperature is not greater than 350xc2x0 C., especially not greater than 330xc2x0 C.\nThe process may be carried out at atmospheric pressure although superatmospheric pressure, say up to about 30 bar is preferred.\nThe contact time is preferably in the range from about 0.1 seconds to about 10 seconds, preferably in the range from about 0.5 seconds to about 5 seconds at atmospheric pressure.\nAs described previously, hydrofluorination of trichloroethylene via the intermediate 1-chloro-2,2,2-trifluoroethane is used for the production of 1,1,1,2-tetrafluoroethane. If desired the process of the present invention may be combined with processes for the production of 1,1,1,2-tetrafluoroethane based on trichloroethylene/1-chloro-2,2,2-trifluoroethane,\nThus and according to a preferred embodiment of the invention, a hydrochlorofluoroalkane and/or a(hydro)chlorofluoroalkene as hereinbefore defined, for example 1,1-dichloro-2-fluoroethene is fed as a second starting material to processes for the production of 1,1,1,2-tetrafluoroethane employing trichloroethylene as the starting material.\nThe co-feeding of trichloroethylene and the second starting material such as 1,1-dichloro-2-fluoroethene may be effected in the processes described in our published European Patent Applications Nos 0 449 617, and 0 449 614, the contents of which are incorporated herein by reference.\nThe invention is illustrated but not limited by the following Examples.\nThe 1,1-dichloro-2-fluoroethene (HCFC 1121a) used in the Examples was synthesised via the ethanolic potassium hydroxide dehydrochlorination of trichlorofluoroethane (HCFC 131) as described in J. Am. Chem. Soc., 1936, 58, 402. The resulting crude product was washed with water, dried with magnesium sulphate and then fractionally distilled with the fraction boiling between 32xc2x0 C. and 34xc2x0 C. being collected; this fraction was analysed and found to be pure HCFC 1121a.\nThe 1,1-dichloro-1,2-difluoroethane (HCFC 132c) used in the Examples was synthesised by oxyfluorination of 1,1-dichloroethane (vinylidene chloride) using lead (IV) oxide in anhydrous hydrogen fluoride as described in J. Am. Chem. Soc., 1945, 67, 1639. The reaction was carried out in a Hastalloy C autoclave and yielded a considerable amount of the co-product 1,1-dichloro-1-fluoroethane (HCFC 141b). The resulting reaction mixture was fractionally distilled and a fraction comprising 60% HCFC 141b and 40% HCFC 132c by weight was collected; it was not feasible to separate the HCFC 132c from the HCFC 141b by distillation."}
{"text": "Data storage devices such as disk drives comprise a disk and a head connected to a distal end of an actuator arm which is rotated about a pivot by a voice coil motor (VCM) to position the head radially over the disk. The disk comprises a plurality of radially spaced, concentric tracks for recording user data sectors and servo sectors. The servo sectors comprise head positioning information (e.g., a track address) which is read by the head and processed by a servo control system to control the actuator arm as it seeks from track to track.\nData is typically written to the disk by modulating a write current in an inductive coil (write coil) to record magnetic transitions onto the disk surface in a process referred to as saturation recording. During read-back, the magnetic transitions are sensed by a read element (e.g., a magneto-resistive element) and the resulting read signal demodulated by a suitable read channel. Heat assisted magnetic recording (HAMR) is a recent development that improves the quality of written data by heating the disk surface during write operations in order to decrease the coercivity of the magnetic medium, thereby enabling the magnetic field generated by the write coil to more readily magnetize the disk surface. Any suitable technique may be employed to heat the surface of the disk in HAMR recording, such as by fabricating a laser diode and a near field transducer (NFT) with other write components of the head. Microwave assisted magnetic recording (MAMR) is also a recent development that improves the quality of written data by using a spin torque oscillator (STO) to apply a high frequency auxiliary magnetic field to the media close to the resonant frequency of the magnetic grains, thereby enabling the magnetic field generated by the write coil to more readily magnetize the disk surface. Since the quality of the write/read signal depends on the fly height of the head, conventional heads may also comprise an actuator for controlling the fly height. Any suitable fly height actuator (FHA) may be employed, such as a heater which controls fly height through thermal expansion, or a piezoelectric (PZT) actuator."}
{"text": "U.S. Pat. No. 5,100,889 to Misra et al discloses 7-oxabicycloheptyl substituted heterocyclic amide prostaglandin analogs which are thromboxane A.sub.2 (TXA.sub.2) receptor antagonists or combined thromboxane A.sub.2 receptor antagonist/thromboxane synthetase inhibitors useful, for example, in the treatment of thrombotic and/or vasospastic diseases, and have good duration of action. Examples of compounds disclosed in Misra et al have the structural formula I ##STR1## and including all stereoisomers thereof, wherein m is 1, 2 or 3; n is 0, 1, 2, 3 or 4;\nR.sup.1 is hydrogen, lower alkyl, aralkyl, aryl, cycloalkyl, cycloalkylalkyl, or amide ##STR2## wherein t is 1 to 12 and R.sub.a is lower alkyl, aryl, cycloalkyl, or cycloalkylalkyl); PA1 R.sup.2 is hydrogen, lower alkyl, aryl, or aralkyl; or R.sup.1 and R.sup.2 together with the nitrogen to which they are linked may form a 5- to 8-membered ring. PA1 R.sup.9 is H, OH or lower alkyl, preferably OH or H. PA1 R.sup.2 is hydrogen, lower alkyl, aryl, or aralkyl; or R.sup.1 and R.sup.2 together with the nitrogen to which they are linked may form a 5- to 8-membered ring. PA1 R.sup.1 is preferably lower alkyl such as n-pentyl, aryl such as phenyl, halophenyl such as 4-chlorophenyl, or cyclohexylalkyl, such as cyclohexylbutyl. PA1 R.sup.2 is preferably H or phenyl. PA1 R.sup.8 is H, aryl or lower alkyl, preferably H and PA1 R.sup.9 is H, OH or lower alkyl, preferably OH or H.\nMisra et al disclose that these compounds may be prepared by transmetallating bromophenylalkyl B ##STR3## by treatment with t-C.sub.4 H.sub.9 Li or n-C.sub.4 H.sub.9 Li or subjecting B to a Grignard reaction by treatment with Mg, and then condensing with the perhydro benzopyran-3-ol derivative or the perhydro benzofuran-1-ol derivative C ##STR4## to form the condensed 7-oxabicycloheptane alcohol compound of the structure Z ##STR5## and then subjecting the condensed compound to hydrogenolysis to form the following alcohol D ##STR6## Where Pro is thexyldimethylsilyl or t-butyldimethylsilyl, the alcohol is acetylated and the silyl protecting group of the so-formed acetate is removed to form the following acetate: ##STR7## which is treated with a protecting compound and the acetate is removed by treatment with aqueous hydroxide or excess methyllithium to form the following alcohol: ##STR8## (where Pro is t-butyldiphenylsilyl). The protected alcohol is subjected to a Jones oxidation to form the following acid: ##STR9##\nThe so-formed carboxylic acid intermediate is then employed to make the final compound.\nIn a more preferred procedure, Misra et al disclose protecting the alcohol function of alcohol Z to form the protected alcohol ##STR10## subjecting the protected alcohol to a Jones oxidation and esterification to form the ester ##STR11## which is made to undergo hydrogenolysis and subsequent removal of the acetate protecting group by transesterification to afford the alcohol ##STR12## which is subjected to a Jones oxidation to form the carboxylic acid intermediate II ##STR13##\nIn an alternative procedure where n is 1, the above carboxylic acid intermediate II is formed by treating D' with acetic anhydride and removing the protecting group to form the acetate alcohol ##STR14## which is made to undergo a Dess-Martin oxidation to form the aldehyde ##STR15## The above aldehyde is oxidized and esterified to the corresponding acetate ester, deprotected, and subjected to a Jones oxidation to form carboxylic acid II where n is 1."}
{"text": "1. Field of the Invention\nThis invention relates to flying or floating craft and to devices which utilize a flowing fluid or movement of a body in fluid to create a rotation and/or lift or thrust.\n2. Description of the Prior Art\nHelicopters were commonly provided with an autorotation arrangemnent. However, the descending speed of the helicopter in autorotation is rather high. The \"descent-speed slow down\"-maneuver is a delicate procedure which often fails and the helicopter crashes. The trend in development of helicopters is presently to reduce the weights. Also the weights of the rotors becomes reduced. That reduces as a side effect the time of the landing maneuver in autorotation and thereby increases the change of failure of the landing maneuver, whereby the likeliness of a crash increases, when the engine of the helicopter fails.\nIn motor powered gliders it was custom in the former art to swing a motordriven propeller into and out of the craft, when a change-over from motorless gliding to powered gliding and vice versa was desired. The swing out of the propeller resulted in a negative influence on the flight path of the glider. The sudden appearance and disappearance of the drag of the propeller brought a negative influence on the speed and stability of the glider in the surrounding air."}
{"text": "Connectors became indispensable parts used for connections in wiring of electrical and electronic equipment used in cars and other applications. Generally such connectors are held together due to the friction between connected contacts and housings, however, the shortcoming of the connections with retention based on friction is that they can not withstand vibration and other impacts, which can lead to their separation.\nTherefore, a number of designs was suggested regarding connectors able to withstand vibrations. Among them, there are connectors with latching devices locking in the state of complete connection that prevent their separation due to vibration.\nHowever, since the contacts of the connectors start to form electrical contact before the state of complete connection is achieved, it is possible to make electrical connection without completely joining the housings.\nIn the manufacturing process, the judgment concerning the state of the connection of the housings is made by visual observation of the state of the latching device or by existence of electrical connections between the wires. However, electrical connection often exists even at an incomplete joining of the housings. Visual observation of the latching devices also can produce erroneous judgment, due to the observation angle or a human error. These factors impart reliability of connectors.\nIn recent years, requirements to the reliability of electrical connectors have sharply increased, especially in relation to a wide use of such safety devices as air bag systems and anti-lock brake systems, etc. Therefore, the problem of an easy confirmation of a complete joining of the connector became especially important. A number of designs of connectors were offered making it possible to easily confirm complete joining.\nFor example, in the connector described in Japanese Utility Publication No. (1992)-3419, an electrical detection device was proposed. In this connector, an active latching device (or a passive latching device) of one connector housing produces an open or a closed electrical contact; when the passive latching device (or an active latching device) of the other connector housing comes against the active latching device (or a passive latching device) of the first connector housing, it switches the condition of the electrical contact to opposite, thus making it possible to confirm connection.\nIn the Japanese Patent No. (1985)-1774, one connector housing has a built-in indicating rod movable in the direction of the joining of the housings. When this connector housing is joined with the matching connector housing, a part of that connector housing pushes the indicating rod which sticks out of the housing by the length required for the connection.\nHowever, the above mentioned connectors equipped with opening and closing electrical contacts are very complicated and expensive to produce. They also require a special detecting circuit, thus making the price of the entire system prohibitively high.\nThe disadvantage of the connectors equipped with the indicating rod is the fact that the length of the exposed portion changes gradually, which may easily end in an erroneous human judgment regarding the state of connection. It also can be used in applications where the connector is out of the vision field of the assembly worker.\nThis invention was made with the purpose to offer low-cost and simple design connectors in which the above disadvantages are corrected, making it possible to make judgment on the state of connection based on direct visual observation or by touch."}
{"text": "1. Field of the Invention\nThis invention relates to an electrical energy saving device for motors, which comprises a principal circuit system operating with a control circuit system for automatically adjusting the magnitude of the input power of the motor with respect to the different load states of the motor.\n2. Related Art\nConventionally, all motor control systems have not been provided with a specially designed triggering device for controlling the conducting condition of the motor. Therefore, when the motor is in an idle running state or the load thereof is light, the input current drawn by the motor is still kept at 50-70% of the full-load current. This mode of operation not only consumes too much energy and reduces the power factor of the system, but also decreases the life span of the motor because of the high temperature suffered by the motor when running at this current level.\nAlthough control elements, such as a SCR or a TRIAC, are used in conventional motor control circuits for triggering the running state of the motor, owing to the fact that the magnitude of the starting current is usually 7-8 times of the zero-load current, the specifications of such control elements must be augmented, which causes a significant increase in the manufacturing costs. On the other hand, a timing device is also used in conventional control circuit as a starting and overload protecting device, but it has to be matched with an electromagnetic switch, which increases substantially the volume of the control device as well as the manufacturing costs. Besides, the timing device, which achieves an overload protection only during the starting period, provides no protection during the running state of the motor. Furthermore, conventional motor control circuit systems cannot be installed at the front end of the power source switch, and are very inconvenient in utilization, replacement and adjustment operations."}
{"text": "Nanoparticles have attracted a great deal of attention in fields such as catalysis, magnetics, and optics, due to their extraordinary properties. In particular, the exploration of new frontiers in catalysis has been intimately related to developments of well-defined nanomaterials. Over the past 20 years, an explosion of interest in nanomaterials has greatly promoted fundamental understanding of catalysis. Effects of particle size, crystal plane orientations, and surface defects on catalytic performances have all been investigated on single component catalysts, partially enabled by the realization of nanomaterials with well-defined morphologies.\nRecently, interest has shifted toward bimetallic catalyst systems, due to their potential to enhance catalytic activity or even create bi-functional surfaces capable of propagating technically challenging chemical reactions. For bimetallic catalysts, outstanding catalytic performance has been reported for various reactions, including alcohol oxidation and oxygen reduction. These successes have triggered intense interest in the preparation of bimetallic catalysts with controlled morphologies and sizes. Among the as-obtained bimetallic nanoparticles, those with core-shell structures are especially interesting. The core-shell nanoparticles have shown superior catalytic activity and/or selectivity in many reactions, which could be partially attributed to high index facets on the surface. Another benefit of the core-shell architecture is a reduced catalyst cost by minimizing the usage of the expensive active component.\nTransition metals are often categorized as being noble and non-noble, in general the morphology and size control of noble-metal nanoparticles are relatively well understood. Despite the diverse spectrum of synthesized noble-metal nanoparticles, synthesis procedures are often similar, with most incorporating the use of capping agents, for instance poly (vinyl pyrrolidone) PVP. Mechanistically it is understood that the capping agents control the morphology and size of the nanoparticles by binding with the noble metals through carbonyl or amino groups.\nIn contrast to the case for noble-metal nanoparticles, non-noble metal nanoparticles are notorious for the difficulties associated with controlling the morphology and size. The binding energies of the capping agents on transition metal surfaces are different (lower in most cases) from that on noble metal surfaces. The disparity can be related to differences in the d-band structure, since the d-orbitals of non-noble metals are filled with fewer electrons than that of noble metals. Filled d-orbital are required to interact with the electrophilic group, carbonyl in PVP, and the electron deficient d-orbitals of non-noble metals thus has lower adsorption energy to PVP. In the literature, to maintain tight control of the morphology and size, two steps are needed to prepare core-shell structures, first forming the core of the nanoparticles and then coating the core with the other component as the shell.\nThough two step can synthesis procedures produce core-shell particles with narrow size distribution, it results in a mono-metallic surface rather than presenting a solid solution of two metals on the surface. The mono-metallic shell is good for reactions with single-active-site mechanism, but not for those with double-active-site mechanisms. Furthermore, the two-step synthesis is prohibiting large scale production due to the complicated procedure, and delicate centrifuging/re-dispersion."}
{"text": "Sunlight is the most abundant renewable energy source available on Earth, supplying over 1.5×1022J (15,000 Exajoules) of energy to the Earth's surface daily. This amount of energy is 10,000 times greater than the 1.3EJ of energy consumed daily by the Earth's population. However, the capital costs involved with harvesting solar energy make supplying energy to the electrical grid using photovoltaic (“PV”) generation unfavorable compared to existing standard technology, e.g., coal, nuclear, etc. Thus, substantial improvements are necessary to justify a shift from existing power generation technologies to photovoltaic-generated energy.\nThe main challenge in the photovoltaic industry is making solar cells more cost effective. Thus, it is desirable to produce less expensive cells with higher efficiencies. To date, the National Renewable Energy Lab (NREL) is believed to have achieved a record efficiency of 40.7% (AM1.5D, low AOD, 240 suns, 25° C.) for metamorphic or lattice mismatched and 40.1% (AM1.5D, low AOD, 135 suns, 25° C.) for lattice matched three-junction photovoltaic cells such as GaInP/GaInAs/Ge cells. However, this is well below the desired efficiency of 60% or higher believed necessary to make PV technology a feasible substitute for existing power generation technologies, so additional improvements are necessary."}
{"text": "Today's computer aided design (CAD) systems are used for a variety of applications, such as mechanical or electronic engineering. A common characteristic of all of these applications is that parts, whether mechanical, electronic or of whatever kind, are designed in interactive mode. That is, the part--as far as already designed--is displayed on a screen (e.g. as two-dimensional, or as perspective view) such as a CRT (cathode ray tube) or an LCD (liquid crystal display). The user enters commands via appropriate input means, preferably a computer mouse, a graphics tablet or a light pen, in order to complement or modify the existing structure. When the editing process is finished, i.e., the part is defined, it may be plotted or otherwise be reproduced. It is also possible to generate punch or magnetic tapes which then directly control a numerically controlled machine tool, in order to manufacture a physical representation of the object.\nIt is understood that it is quite easy to generate or modify parts (or their contour) of simple geometry, such as straight lines, circles, arcs, cylinders, spheres and the like, as these can be expressed and calculated in simple mathematical terms and equations. However, the task becomes more complicated for surfaces and boundary lines which cannot easily be expressed in analytical mathematical terms, but should still intentionally be smooth curves. A well-known example of this kind is the body of a vessel.\nTherefore, in a variety of applications, it has become common practice to use a low-degree polynomial function (in some cases, even functions of other types, such as exponential or trigonometric functions) to describe non-regular contours of a physical object. A well-known function of this type is the B-Spline which approximates a non-regular contour. The contour or curve is derived from a control polygon (a piecewise linear function associated with the B-Spline) by means of a matrix multiplication of the control points (corner points of the control polygon) with the so-called basic functions of the B-Spline.\nThis will be illustrated now in more detailed manner. As an example, the B-Spline approximation has been chosen. However, it has to be noted that this selection has only been made for convenience; the present invention may be applied to other approximation techniques (e.g., the Bezier or Hermite form) as well.\nIn general, a B-Spline can be denoted in parametric terms as follows: ##EQU1## wherein Q(t) denotes the curve (contour of the B-Spline)), the B.sub.i (t) are the basic functions of the B-Spline (rational or non-rational), and the P.sub.i are the control points of the control polygon. As the Spline curve is usually a graph in two- or three-dimensional space, the above equation is a vector equation; the vectors are denoted as bold characters herein. \"i\" is an index running over all control points and basic functions.\nIn a more general case, the geometric object is a surface in space, approximated by a B-Spline surface of the notation ##EQU2## with u,v as (independent) parameters, the basic functions M.sub.i (u) and N.sub.i (v), the control points P.sub.ij and the indices \"i\" and \"j\".\nThere are basically two ways of defining a B-Spline curve or surface. The first is to define the control points P.sub.i of the control polygon; the B-Spline curve. is then calculated from the control points according to the above equation. This way requires only few programming efforts and computing power; however, as the curve or surface is generated indirectly, the user does not have full control over the final geometry, such that several iterations may be necessary until the desired curve or surface is designed.\nThe second way is based on the definition of interpolation points which are located on the function (B Spline curve or surface). The B-Spline is then generated by interpolation of several B-Spline profiles. It is understood that this way is more convenient to the user, but requires higher programming effort and more computational power.\nA common problem in either of these cases is the modification of already defined (\"existing\") B-Spline curves or surfaces (this applies also to other functions representing a geometric object). The prior art provides two solutions to this problem:\na) In case the Spline curve or surface was generated via the control points of the associated control polygon, single control points can be moved. The Spline is then rebuilt on the basis of the modified control polygon. PA1 This kind of modification operates in quite indirect manner, i.e., the user cannot exactly predict what will happen with the curve or surface when a control point is shifted. In order words, the modified Spline will seldom meet his expectations, such that further iterations are necessary. Needless to say that this process is very time-consuming. PA1 b) In case the Spline curve or surface was generated via interpolation of multiple Splines, it is possible to access interpolation points, or Spline profiles, directly and shift them to another location. PA1 This process provides better control for the user, i.e., the modified Spline will meet his expectations better than in case a). However, this technique has other disadvantages. PA1 The most serious drawback is that the modification of the Spline is not locally restricted. I.e., a modification will result in a modification of the whole Spline, instead of only a local region thereof. This effect is usually not desired. Likewise, the modified Spline will oscillate around the interpolation points. PA1 Further, access is only possible to the interpolation points, but not to other points of the curve or surface. This limits the user's choices to modify the Spline. PA1 identifying at least one point of origin on the geometric object, PA1 identifying a target point for replacing the point of origin, PA1 transforming the shift of the point of origin to the target point into an at least local or regional shift or modification of the at least piecewise polynomial function, PA1 defining a modified geometric object as a function of the modified piecewise polynomial function. PA1 The selected point on the unmodified, and/or the modified, geometric object (or the Splines representing them); and/or PA1 the Spline functions themselves; and/or PA1 the parameter value of the selected point(s) of the unmodified, and/or the modified, geometric object (or the Splines representing them); and/or PA1 form or contour parameters arbitrarily selectable by the user, such form or contour-parameters defining the shape, and the local extent, of Spline modification.\nThe disadvantages and effects of the prior art techniques for Spline modification will be illustrated, and discussed further, in the detailed description.\nConsequently, there is a need for an improved modification technique which avoids the disadvantages encountered with the prior art modification techniques, either partially or completely."}
{"text": "1. Field of the Invention\nThis invention generally relates to the field of electric power generation and electric power conversion, and in particular to a system for maintaining the voltage level of a power supply bus coupled to a load.\n2. Description of the Related Art\nThe need for mobile commercial-grade alternating current (AC) power comes in many applications and cuts across many industries. There is an emerging need for a mobile power generation system that can satisfy the demands of power hungry machinery, like heavy construction tools, while still being capable of delivering well-regulated power to sensitive equipment, like computers. Moreover, not only should the mobile system have the capacity to generate such power while operating at a stationary site, it should also provide the same type of power while the system is in transit. For example, a mobile power generation system should be versatile enough to support the power requirements of emergency response vehicles rendering life-saving services while in transit, as well as satisfy the power demands of sensitive intelligence/news-gathering computer and telecommunications equipment in isolated regions of the world.\nThere is also a need to provide this same AC power during occasions when the running of the vehicle\"\"s engine is not desirable. Late night work in residential areas, where engine noise is an issue, is an example of one such occasion. Generation of power without the need to run an engine is also desirable to reduce air pollution caused by engine exhaust.\nOne solution for this need is to derive the power from storage batteries, converting the energy from a direct current (DC) source to AC. A device that performs such a function is called an xe2x80x9cinverterxe2x80x9d. Inverters are in the prior art and are commercially available. The conventional inverters typically depend on either the charging system of the vehicle to restore charge to the batteries once they are drained, and/or by a plug-in battery charger. It is common for a vehicle with an inverter to be plugged into standard utility power overnight to xe2x80x9ctop offxe2x80x9d the batteries then depend on the vehicle\"\"s alternator to replenish charge level during operation. It is also common for a vehicle to be outfitted with a high output alternator to provide additional charging capacity.\nThe conventional inverters suffer from various disadvantages. For example, in some of the conventional inverters, the power demands from the inverter exceed by a large margin the ability of an alternator to provide charge. In such a case, there comes a point when the batteries are completely drained, and because of the alternator\"\"s inability to keep up alone, inverters shut down and work requiring AC power ceases.\nAnother disadvantage of the conventional inverter system is that a vehicle alternator produces its least amount of charging current at low engine RPM\"\"s. When the batteries are drained, the typical response is to start the vehicle\"\"s engine to recharge the batteries with the alternator. At idle RPM the alternator is limited in capacity and will tend to overheat, reducing alternator life and potentially creating a fire danger.\nA further disadvantage of the conventional inverter system is that the inverter typically continues to pull power off of the batteries until the batteries are very deeply discharge. Deep discharging of batteries may greatly limit the life of the batteries.\nYet another disadvantage of the conventional inverter system is that the typical automotive alternator recharges based on voltage level only, not battery temperature, which in some cases limits the amount and rate of battery charging and reducing battery life.\nYet another disadvantage of the conventional inverter system is it typically lacks the capability of providing very high power, short duration demands, such as those that occur when starting certain AC electric motors.\nAccording to one aspect of the invention, a power generation system is provided for supplying electrical power to a load from at least one of a generator and a battery. The power generation system includes a power bus coupled to the generator and a bi-directional conversion unit coupled between the power bus and the battery. The bi-directional conversion unit is capable of transitioning between a first direction wherein electrical power flows from the power bus to the battery and a second direction wherein electrical power flows from the battery to the power bus.\nIn one aspect of one embodiment, the present invention provides a way to continually supply AC power, even when the generator is no longer producing electrical power, by converting DC power from a battery source into AC power.\nIn another aspect of one embodiment, the present invention provides a system that transitions between a battery power operation to a generator power operation without AC power interruption.\nIn another aspect of one embodiment, the present invention provides a smooth transition from a mode where a battery source is supplying power to AC load to another mode where a generator source is supplying power to the AC load.\nIn another aspect of one embodiment, the present invention is capable of simultaneously supplying DC and AC power without interruption.\nIn yet another aspect of one embodiment, the present invention provides means for supplying power from both a generator source and a battery source during times of high power demand.\nIn yet another aspect of one embodiment, the present invention provides an alarm to an operator that the charge levels of the batteries are getting low, providing time for the operator to start the engine and begin recharging the batteries.\nIn yet another aspect of one embodiment, the present invention monitors battery temperature and adjusts the charging rate accordingly, enhancing battery life and maximizing the amount of useful charge stored in the battery.\nIn yet another aspect of one embodiment, when the electrical power generated by the generator exceeds the AC load demand, the excess power is used to recharge a battery source."}
{"text": "This invention relates to the alleviation of industrial air and water pollution and particularly to the removal of sulfur pollutants from industrial exhaust gases.\nIndustrial gases such as coke oven gas, natural gas and various artificially-produced fuel gases are used either by industrial plants to make useful products or burned in suitable combustion apparatus to produce heat. These gases are composed of varying mixtures of hydrogen, carbon monoxide, various aliphatic and aromatic hydrocarbons, hydrogen sulfide, hydrogen cyanide, carbonyl sulfide and other combustibles. The presence of sulfur compounds in such industrial gases is particularly undesirable because of possible corrosion of intermediate gas transmission lines and other apparatus by the gases, possible contamination of chemical substances made from the gases, and possible discharge of undesirable concentrations of sulfur oxides to the atmosphere during combustion of the gases.\nIn the past such industrial gases have often been treated by passing them through absorption-desorption processes of various types. These absorption-desorption processes give off so-called foul gases which are treated to recover the sulfur present in the gas and thus prevent its discharge to the atmosphere.\nTypical absorption-desorption processes are the hot potassium carbonate process, the vacuum carbonate process, the amine processes, especially those using mono-di-, and triethanolamine, and various other processes using organic solvents. The alkanolamine processes and particularly the diethanolamine and monoethanolamine processes have proven to be especially attractive from an industrial standpoint due to their attractive economics and relatively trouble-free operation. The monoethanolamine processes in particular have proven to be very convenient and efficient in removing hydrogen sulfide and other sulfur compounds from sulfur-containing gas streams. Monoethanolamine solutions can easily remove substantially the entire sulfur content from industrial gases so that the gas leaving the monoethanolamine absorber contains no more than 10 grains of sulfur per 100 standard cubic feet of gas exhausted, a very small amount.\nThese absorption-desorption type processes, while effective to reduce the sulfur content of treated industrial gas to a very low level, do recover the acid gases in a more concentrated form. The recovered \"foul\" gases have to be treated in turn to remove their sulfur content in some satisfactory manner. Very frequently the foul gases from the absorption step have been used to produce elemental sulfur by some variation of the so-called Claus process. In this process a portion of the sulfur removed, usually in the form of hydrogen sulfide, is oxidized to sulfur dioxide and the sulfur dioxide and remaining hydrogen sulfide are then reacted in a catalytic converter to form elemental sulfur and water.\nThere are a number of industrial variations of the basic Claus process in which either an initial portion of the hydrogen sulfide is oxidized to sulfur dioxide or a portion of the final elemental sulfur product is subsequently oxidized to sulfur dioxide for use in the Claus reaction. The Claus reaction may be conducted either in the gas phase by mixing the two gases in appropriate apparatus, in which case the process is usually referred to either as the Claus process or a gas phase sulfur recovery process, or the reaction may be conducted in the liquid phase. The liquid phase process is conducted by dissolving one or both of the component sulfur base reactant compounds in a liquid and then either passing the other reactant compound through the liquid or bringing together two parts of liquid in which the reactants are dissolved. Such liquid reaction medium processes, which are frequently referred to as liquid phase sulfur recovery processes, are usually run at lower temperatures than the more usual Claus type processes and have certain other advantages as well.\nThe Claus process and the other related processes for the recovery of elemental sulfur such as the liquid reaction medium processes are fairly efficient, but have the disadvantage that there is invariably a residue of gas known as the tail gas in which either sulfur dioxide or hydrogen sulfide or frequently both remain. This tail gas must be disposed of in some manner and is usually at this point discharged to the atmosphere.\nWhile the total amount of residual sulfur compounds contained in the tail gas and discharged to the atmosphere is much reduced from the concentrations of sulfur in the original gas treated, there is still, due to inherent inefficiencies of the system, a residual amount of sulfur in the remaining tail gas which may be objectionable. The amount of this remaining sulfur can be decreased by subsequent processing, for example, by the use of several Claus type reactors in series, but due to the small amount of remaining sulfur components in the final tail gas any further processing becomes more and more inefficient and expensive and there is a final minimum of sulfur which is impossible to remove.\nOne fairly simple expedient for final treatment has been to oxidize all the remaining sulfur compounds to sulfur dioxide and then wash the sulfur dioxide out of the gas with a simple water wash system. The wash water can then be discarded, or used if it is concentrated enough to make sulfuric acid. However, the amount of sulfur dioxide dissolved in the water is insufficient for really effective use as a source of sulfur, yet it is undesirable to waste the sulfur values by discarding the wash water.\nThe Claus process in particular exhibits a fairly poor recovery of sulfur based upon the amount of sulfur in the original gas and it is customary to use three or even four Claus reactors in series in order to effect recovery of more than 85 to 95% or occasionally as much as 97% of the sulfur removed from the original gas.\nSome of the many liquid reaction phase processes are, on the other hand, significantly more efficient so far as the percentage of the original sulfur content which is removed from the gas as elemental sulfur is concerned. However, there is always some small amount of sulfur left in the final tail gas which it is substantially impossible to remove. One of the difficulties is that even with the best control and mixing of the two gases H.sub.2 S and SO.sub.2, it is inevitable that one or the other of the two reactant gases will be present in some degree of excess and such excess will pass unreacted through the process. It is simply impossible under industrial conditions to exactly meter the two gases together in an exact stoichiometric relationship. Usually, of course, other reacting conditions will also not be completely favorable for completion of the reaction so that a quantity of one gas or the other, and very often both, is left unreacted. Sometimes it is decided beforehand which gas is most desirable or least detrimental to have in excess and this gas is then deliberately supplied in a slight excess in order to assure that substantially none of the other gas remains at the end of the process. Usually SO.sub.2 will be preferred in the tail gas rather than H.sub.2 S because SO.sub.2 is less objectionable to the senses than H.sub.2 S in similar concentrations.\nThe relative amounts of hydrogen sulfide, carbon dioxide, water vapor and extraneous hydrocarbons in the foul or acid gas fed to Claus reactors have a pronounced influence on the recovery efficiency for a given number of reactors. As pointed out above there is a limit to the number of reactors which can practically be used to reduce the sulfur compounds in the tail gas. There is also a limit to the preciseness with which an operation can be run in order to keep the reaction exactly in balance and prevent decreases in the conversion of the sulfur compound gases to elemental sulfur due to variations from precise stoichiometric ratios of the reacting gases. The reversible reaction EQU 2H.sub.2 S + SO.sub.2 .revreaction. 3S + 2H.sub.2 O (1)\nof the Claus process cannot, furthermore, be completed to the right of the equation because of limitations of the thermodynamic equilibrium at temperatures above 250.degree. C. at which the Claus reaction is run. Sufficient water and sulfur vapor is always present to limit the desired reaction. Carbon dioxide and additional water vapor in the feed gas are also diluents which shift the equilibrium adversely. Exact air/acid gas ratio control is also never continuously achieved even with the most sophisticated instrumentation, because of uncontrollable variations in the feed rate and composition of the gases. Hydrocarbons in the feed also affect the plant efficiently by increasing undesirable side reaction products such as COS and CS.sub.2, which are difficult to convert to elemental sulfur. The tail gas from a single Claus reactor may as a result contain as much as 10% of the sulfur originally removed by the absorption system from the fuel gas.\nA number of processes have been proposed as \"clean-up\" processes for further treatment of Claus unit tail gas. Several of these depend upon treatment of the tail gas so that the residual sulfur values occur as hydrogen sulfide which is then converted to sulfur in a so-called Stretford unit. There are a number of other proposals for improved clean-up of the tail gas including the use of improved catalysts in the Claus reactor, wet scrubbing, reaction with ammonia, high temperature sulfur dioxide removal, concentration of sulfur dioxide by absorption, catalytic sulfuric acid production, absorption on carbon and absorption-desorption type chemical removal. Some of these proposals are applied to the tail gas after incineration to change all the sulfur values to sulfur dioxide. While some are fairly efficient, at least in the laboratory, in removing the sulfur components, many have proved impractical under industrial conditions and none is completely efficient in removing sulfur compounds.\nOne notable yet fairly typical example of a Claus-type process in which residual H.sub.2 S is oxidized to SO.sub.2 is shown in U.S. Pat. No. 1,915,364 issued June 27, 1933 to J. W. Harrel. Harrel separates a gas having a high content of H.sub.2 S into two portions. One portion is sent to a burner or furnace where it is burned with O.sub.2 to SO.sub.2. The SO.sub.2 is absorbed into a water solution which is then passed concurrently or countercurrently to the remainder of the H.sub.2 S gas in a reaction tower where the Claus reaction converts H.sub.2 S and SO.sub.2 to elemental sulfur and water. In order to eliminate side reactions and possible loss of sulfur an excess of H.sub.2 S is maintained in the reactor. This excess results in residual H.sub.2 S gas leaving the reactor with the solution, which H.sub.2 S is then separated from the solution and returned or recycled back to the burner or furnace where it is burned or oxidized with oxygen to form additional SO.sub.2 required for the Claus reaction. Some of this SO.sub.2 eventually escapes through a vent on the top of the absorber where the SO.sub.2 is absorbed by water, but no H.sub.2 S leaves the apparatus and pollution of the air with H.sub.2 S is thus said to be completely eliminated. The elemental sulfur formed in the Claus reaction is separated from the water suspension by any suitable method such as filtering or the like.\nHarrel is fairly typical of prior processes in which excess H.sub.2 S is oxidized subsequent to the claus reaction. Many later processes have, however, attempted after oxidizing the H.sub.2 S to remove the resulting SO.sub.2 from the tail gas by the use of additional removal apparatus, none of which has proved to be really practical or efficient under industrial conditions. A number of processes for the treatment of tail gas from sulfur recovery plants for the minimization of sulfur emissions are summarized in the following articles.\n(a) \"Reduce Claus Sulfur Emissions\" by C. B. Barry Hydrocarbon Processing April 1972 pages 102-106 PA1 (b) \"Production of Clean Energy and Sulfur Recovery\" by Mikio Harima Chemical Economy and Engineering Review, Vol. 6 No. 8 (No. 76), August 1974, pages 13-21 et seq.\nRecycling per se in connection with sulfur recovery processes is, of course, not in itself new as illustrated by the Harrel patent noted above, and a large number of processes have been developed which specifically make use of the broad principle of recycling in order to increase the recovery of the sulfur compounds from a gas. For example, normal Claus reactor tail gas containing both H.sub.2 S and SO.sub.2 has been oxidized to convert all residual H.sub.2 S to SO.sub.2 and the SO.sub.2 has then been recycled back to the Claus reactor to replace a portion of the SO.sub.2 used in the reactor. In some proposals the SO.sub.2 has been absorbed from the tail gas into lime or the like and then regenerated from the lime and recycled into the Claus reactor.\nSeveral processes have been developed in which H.sub.2 S and SO.sub.2 are reacted together in a liquid reactor medium of some suitable composition. The liquid reaction medium may be renewed at intervals by stripping gases and volatilizable components including, for example, H.sub.2 S which is then recycled to the primary reactor. Occasionally an ammonium salt solution such as an ammonium sulfite solution has been used as the reaction or absorption solution and in these cases excess H.sub.2 S or SO.sub.2 passing from the solution may be recycled back to reconstitute the absorption solution. It has also been proposed to recycle the entire tail gas stream containing both SO.sub.2 and H.sub.2 S from a Claus reactor back to an original coke oven gas stream to react with the ammonia in the coke gas. More recently it has also been proposed to use a so-called molecular sieve to reversibly adsorb H.sub.2 S from a tail gas derived from a Claus reactor and recycle it back to an absorption step.\nWhile prior workers have, therefore, used the principle of recycling in various ways in connector with sulfur removal systems for the desulfurization of industrial gas, such prior systems have not been completely successful in eliminating exhaust of sulfur components to the environment and have in many cases been expensive to build or uneconomical to operate. Furthermore, while a great many prior workers have made use of the basic principle of recycling various sulfur compounds back to various parts of the processes for retreatment in order to increase the recovery of the sulfur compounds, none has conceived or used the novel arrangement developed by the present inventors."}
{"text": "The invention generally relates to the field of prostheses for the vertebral column, which are intended to correct and stabilize scoliosis or to reposition and stabilize vertebrae in the case of a trauma or a tumor.\nSuch prostheses, called spinal osteosynthesis devices, can be implanted either by the posterior route or by the anterior route. When a posterior route or approach is used, the surgeon is unable to directly access the vertebral bodies. In contrast, the anterior route has in particular the characteristic of permitting direct access to the vertebral bodies, as well as many other advantages which are listed in the document FR-A-2 697 744, which is hereby incorporated by reference in its entirety.\nSuch devices implanted by the anterior route and permitting fusion between two or more adjacent vertebrae can be classified in two categories:                those which use one or more components in the form of rigid plates or the like, the plates being fixed to the different vertebrae in the zone to be treated;        those which use one or two metal rods fixed along the vertebral column by means of suitable connection devices.        \nTo fix these devices on the vertebral column using the known techniques, the surgeon fixes, on each vertebra of the zone in question, a plate-like anchoring member, called an anchoring plate, which serves as an anchoring point for the component or components for regulating the relative positions of these vertebrae, namely the plates and/or the rods.\nGenerally, the anchoring plate is placed in an orientation that is substantially parallel to the vertebral discs, and its lower face must be quite strongly curved inwards to best match the surface of the vertebra in this orientation. In the case where plates are used, these are directly fixed on the anchoring plate. In the case where rods are used, these are generally connected to the anchoring plate via one or more connectors that can each connect a single rod to the anchoring plate (as in FR-A-2 761 256), or two rods simultaneously (as in FR-A-2 697 744).\nIn order for the existing devices to be implanted in the thoracic region, it is necessary for the surgeon to make an incision between two ribs. In some cases the surgeon is forced to section a rib. However, as the available space is relatively small, the surgeon then has few possibilities for fixing another anchoring plate on an adjacent vertebra. Therefore, in order to do this, the surgeon very often has to make a new incision, with all the disadvantages which this means for the patient."}
{"text": "Many social networking services, such as Facebook or the professional social networking service LinkedIn®, make recommendations to their users. These recommendations may pertain to people with whom to connect, articles to read, jobs for which to apply, etc. The quality and relevance of such recommendations may be heavily dependent on the underlying representation of various content items used to generate such recommendations. Examples of content items or objects are a member profile, a job posting, a SlideShare article, a Pulse article, etc.\nToday, the quality of many recommendations suffers from the problem of vocabulary mismatch between different content types. For example, if a member profile of a member of a social networking service (also referred to herein as “SNS”) and a job description use different terminologies to refer to the same underlying concept, the SNS may fail to match the member profile to the job description, and to recommend the respective job to the member.\nTo address this problem, it may be beneficial to an SNS to generate a universal concept graph that includes a unified and standardized set of concept phrases that may be used to generate better recommendations to the members of the SNS."}
{"text": "Mobile computing devices have been developed to increase the functionality that is made available to users in a mobile setting. For example, a user may interact with a mobile phone, tablet computer, or other mobile computing device to check email, surf the web, compose texts, interact with applications, and so on. One challenge that faces developers of mobile computing devices is efficient power management and extension of battery life. For instance, the small form factor of many mobile computing devices may compel designs in which power connections are kept relatively small in size. Accordingly, developers may be further concerned with ensuring that adapters designed for use with such small form power connections are safe and supply the proper amount of power to the device."}
{"text": "Photodiodes operating in the spectral range with a wavelength of up to approximately 1.1 μm are usually produced from silicon. Photodiodes composed of germanium or InGaAs are sensitive in the infrared range. By using further compound semiconductor materials, for example II-VI semiconductors or else organic semiconductor materials, further spectral ranges can furthermore be covered.\nThe p-n junction in the semiconductor material is produced by using regions having differing doping which provide for a change of charge carrier type. In the region of the p-n junction, a depletion layer, also called space charge zone, depleted of charge carriers is formed on account of diffusion and recombination processes, in which depletion layer an electric field prevails which provides for a separation of the charge carriers in the space charge zone. On account of the continuous thermal generation of charge carriers in the space charge zone and in adjacent regions of the neutral zone from which the charge carriers migrate into the space charge zone as a result of diffusion, a small current, also called reverse current, flows permanently via the space charge zone.\nIf electromagnetic radiation impinges on the space charge zone or, respectively, the adjacent diffusion region, further electron-hole pairs are generated which, in addition to the permanent reverse current, provide for a photocurrent generated by the electromagnetic radiation. In this case, the photocurrent is proportional to the intensity of the electromagnetic radiation. The electromagnetic radiation has to have a minimum energy in order to generate the charge carriers, the known energy being dependent on the band gap of the semiconductor material used. A limit wavelength above which the sensitivity decreases greatly thus results for each semiconductor material. Furthermore, the penetration depth of the electromagnetic radiation is also greatly dependent on the wavelength, the penetration depth increasing with longer wavelengths. In the case of silicon, for example, electron-hole pairs are generated by using visible light up to a limit wavelength of 1.1 μm, approximately 90% of the incident light radiation being absorbed at a depth of 0.5 μm in the case of blue light, up to depth of 1.1 μm in the case of green light and up to a depth of 8.5 μm in the case of red light.\nIn order to detect the colour information of electromagnetic radiation over a relatively large wavelength range with the aid of photodiodes, as is carried out for example during the detection of colour images with the aid of photodiode arrays in digital cameras, various solutions are known. Thus, for colour resolution, a plurality of photodiodes each coordinated with the spectral range to be detected are interconnected, wherein the individual photodiodes are provided with spectral filters in order to transmit the desired spectral range to the associated photodiode. In this case, the spectral filter can be embodied as an absorption filter for masking out non-relevant radiation and can be used in a Bayer pattern. Furthermore, transmission gratings are known, which are applied over the photodiodes and act as spectral filters in order that only electromagnetic radiation having a specific wavelength is transmitted. Such a construction is described for example by P. B. Catrysser, B. A. Wandell in: Journal of Optical Society of America, Vol. 20, No. 12, December 2003.\nFurthermore, methods are known wherein the electromagnetic radiation is spectrally split with the aid of prisms and fed separately to the assigned photodiodes. Furthermore, there is the possibility, in the case of photodiodes which are sensitive over a wide spectral range, of additionally using discrete filters which are fitted alternately above the photodiode structure in order that the individual spectral ranges are detected sequentially one after another.\nFurthermore, in the case of photodiodes, the space charge zone can be expanded by applying a reverse voltage, in order, with the aid of the larger space charge zone, to improve the effectiveness of the charge carrier detection particularly for long-wave electromagnetic radiation having a large penetration depth. However, the range in which the expansion of the space charge zone can be varied under typical operating voltages for photodiodes is predetermined by the doping concentration and is generally not adjustable over the range of approximately 0.1 to 10 μm which is of interest for visible light.\nFurthermore, individual photodiodes with spectral resolution are also known wherein no additional spectral filter is required. Thus, U.S. Pat. No. 5,965,875 describes a photodiode construction wherein a plurality of p-n junctions are formed at different depths in the semiconductor material, as a result of which a plurality of space charge zones are formed one above another and separately from one another, which can then be used to simultaneously detect different spectral ranges of the incident light radiation. In this case, each p-n junction has a dedicated electrical connection for reading out the photocurrent detected in the respective spectral range.\nWO 2008/068616 furthermore describes a photodiode construction including a large space charge zone, in which, as a result of transverse electric fields being applied, the charge carriers generated at different absorption depths in the space charge zone are separately collected and read out, such that different spectral ranges of the incident light radiation can be resolved. A similar principle is employed in US 2006/0244089, wherein a large-area electrode is applied above a photoactive layer with a large space charge zone, which electrode is insulated from the photoactive layer and can be used to control the extraction region for the charge carriers in a manner dependent on the absorption depth, in order to obtain a spectral resolution of the incident light radiation.\nThe implementation of spectrally sensitive photodiodes, also referred to as colour sensors hereinafter, as clusters composed of a plurality of photodiodes specifically tailored to the spectrum that is respectively to be detected, or as broadband photodiode with additional colour filters is complex in production and can be integrated only with difficulty into standard processes for producing semiconductor circuits, in particular in CMOS circuits. Furthermore, a large area requirement is generally necessary, which limits the image resolution when such colour sensors are used in CMOS cameras, for example.\nMoreover, generally only a discrete colour resolution is possible with the known colour sensors. This also applies, in particular, to photodiodes having multiple p-n junctions in order to provide different absorption depths and hence spectral sensitivities. Furthermore, each p-n junction has to be individually electrically connected in order that the charge carriers generated can be read out separately and the individual colours can thus be resolved.\nThe design of colour sensors as lateral photodiodes having a multiplicity of electrical connections in order, by using corresponding transverse fields, to read out the charge carriers from different absorption depths and thus to obtain a colour resolution demands a large area requirement and is furthermore complex in production. In the case of the photodiode structure proposed in US 2006/0244089, wherein control is obtained with the aid of an insulated surface electrode, there is the difficulty that the electrode material already absorbs electromagnetic radiation prior to penetration into the actual photoactive layer. As a result, particularly in the short-wave spectral range, less light intensity is available for electrical signal formation.\nFor these and other reasons there is a need for the present invention."}
{"text": "Bag computers are designed to hold a pivoting display panel on their front wall so that the computer can be used with little set up and used a mobile environment while a standing, sitting or lying position. While the bag's shoulder strap holds the bag to the operator, he can use two hands to use the computer controls. Thus, it is a wearable computer.\nBag computer have a pivoting display in their front panel. The display can be stored flat against the bag front and then can be pivoted to about 90 degrees for viewing. The display panel may also be pivoted to 180 degrees so the bag can be set upright in the operator's lap while sitting and viewed over the top of the bag thus bringing the display closer to the user and reducing neck strain.\nOne way to make a bag computer is to have a two panel computer component mounted to a matching bag front using an outside computer holder or inside mounting structure and coupling. The two panel computer has a body panel and display panel pivotally connected with the display fixed to the display panel side facing away from the body panel. The two panel computer may further be installed in the bag's matching computer holder with it's hinge near the bag top front so that the display panel, when pivoted into the line of sight of the operator/bag wearer, is as close as possible to the operator's eyes. The hinge connecting the two panels may include an extension which creates a gap between the two panels so that the material of the bag's computer holder can fit between the two panels when they are closed together in storage position. This arrangement was described in U.S. application Ser. No. 11/796,920.\nImprovements to the hinge and hinge area of the two panel computer can; 1) improve the installation and fit of an inside mounted two panel computer to the bag; 2) allow for a larger, stronger hinge which can also be a venue for related accessories; 3) allow for adjustment of the attachment position of the hinge to the computer body. This application provides embodiments for these improvements."}
{"text": "Millions of consumers regularly use product redemption coupons and realize substantial savings as a result. Significant time is spent clipping and sorting coupons, discarding expired coupons and organizing current coupons for use on shopping trips. Conventional coupon distribution results in significant wasted time due to consumers' attempts to manage their coupon use.\nCoupons are delivered to consumers through a variety of media such as, for example, free standing inserts (FSIs), in-store shelf coupon dispensers, check-out coupons (generally issued based on the customer's current purchase), register receipt coupons, in-product coupons, instant peel-off on-product coupons, direct mail coupons, and/or other delivery media. In addition to manufacturers' coupons, consumers use retail store coupons, such as those issued by large retail chains on a weekly basis. Some consumers use coupons on a fairly random basis. These consumers tend not to keep coupons for future use, but will review coupons available just prior to shopping to see if any of them cover products they plan to buy or if there are any for new or improved products of interest.\nMore organized coupon users maintain some form of storage system to keep coupons for future use. These consumers often clip coupons regularly from all available sources, and often have coupon filing systems by product category. They will also review their coupons regularly, discarding unused coupons which have expired.\nFor most consumers, attempts to maintain an organized coupon file often fails. The “bother” and time required to maintained organized coupon files often results in neglect of those files, even though diligent shoppers know that a consistent significant savings is easily achievable using coupons.\nThe notion of issuing product redemption coupons to consumers was an innovative idea to entice consumers to try new products in the hope that, after the first try of a new product at a coupon discounted price, they would become repeat customers at the regular price. Coupons are effective tools used in launching new products. Manufacturers also find coupons can shore up flagging sales, help reduce excess inventory or win back consumers' brand loyalty, and so coupons for existing products have become customary, so much so that today's consumers have come to expect coupons. Often, coupon price incentives significantly reduce brand loyalty, and manufacturers must issue more coupons than desired to maintain market share. Market share also has been impacted by an increase in the number and variety of competing “no-name” store brands. The competitive nature of the retail industry does not allow manufacturers to reduce coupon distribution, and in some market sectors, such as cereals, the majority of purchases are made with coupons.\nAttempts have been made to meet the needs of the coupon industry and the consumer. For example, attempts have included television-based coupon reception system wherein coupon information is transmitted along with program information to a broadcast audience, electronic paperless coupon system where coupon signs are placed near the related item and the discount is automatically applied at the point of sale, kiosk type printer stations located at a retail store linked to the manufacturer(s) in order to obtain specific coupon information, and/or other attempts. While these aforementioned attempts at providing couponing systems are useful in their own right, they fail to provide for a secure and interactive coupon generation system in which the user can request, select, store, manipulate and print coupons as desired, in which user-specific information such as demographic data and data representative of those coupons so requested, selected, printed and actually used may be provided back to the coupon issuer and distributor for more efficient coupon targeting in subsequent coupon issuance and distribution.\nFurther, the aforementioned attempts at providing coupon systems fail to provide coupons or promotions which are customized according to user-specific information, such as user demographic data, including, for example, location data, proximity data, competitor data, time data and/or date data. Thus, prior attempts do not provide electronic systems or methods which allow a marketer of a product to customize a promotion in such a way as to maximize the value of the promotion to both the marketer of the product and a consumer of the product. In particular, the prior attempts do not allow a promotion which is customized according to the proximity of a consumer from a point of sale and which factors into consideration the proximity of the consumer to a competitor or competitors.\nIt is therefore an object of the present invention to provide such a coupon distribution system which overcomes the aforementioned problems and shortcomings of the above-described references.\nIt is an object of the present invention to provide an electronic coupon distribution system which can be easily accessed by masses of consumers by using a readily available personal computer rather than needing to purchase special-purpose equipment.\nIt is a further object of the present invention to provide such an electronic coupon distribution system which allows a user to request transmission of incentive data and select, store, manipulate and print coupons from such incentive data.\nIt is a further object of the present invention to provide such an electronic coupon distribution system which allows the coupon issuing companies to access valuable information directly from the consumer without requiring specific and additional action by the consumer but rather by using the information from the user's personal computer regarding the consumer's selection, printing and actual redemption of coupons, as well as responses to demographic queries posed to the users.\nIt is a further object of the present invention to provide such an electronic coupon distribution system which allows a consumer to generate shopping lists associated with coupons selected and printed, in order to simplify the shopping process and promote the use of product coupons.\nIt is a further object of the present invention to provide such an electronic coupon distribution system which allows for automatic deletion of expired coupons in the user's computer database and the modification of redemption amounts of coupons in the user's database, both of which can be transparent to the user.\nIt is a still further object of the present invention to provide a secure coupon system which generates unique coupons with user-identifying data and allows the printing of a coupon only once, thus eliminating the possibility of fraud by both the consumer and the retailer.\nIt is a still further object of the invention to provide an efficient, low cost, zip-code/lifestyle/lifestage or household targeted coupon distribution system to tailor the incentives to each user.\nIt is a still further object of the invention to provide a system which allows a marketer to customize a promotion, including customizing a promotion based on user demographic information, including user location information, and/or user lifestyle information.\nIt is a still further object of the invention to provide a system which provides incentive data or coupon information to consumers, wherein the incentive data or information is customized to a specific class of consumers, such as consumers within a specified geographic area or consumers within a specified proximity to a retail center or competitor-retail centers.\nIt is yet a further object of the invention to provide a promotion which is customized to a class of consumers based upon the geographic location of the consumers or the proximity of the consumers to a retail center having the product corresponding to the promotion and which factors into consideration the proximity of the consumers to alternative retail centers, including competitor-retail centers.\nIt is yet a further object of the invention to provide coupons which are customized to a specific class of consumers, such as consumers in a specified geographic location, consumers in a specified proximity to a retail center or consumers with certain common demographic or lifestyle information, including coupons which are capable of being associated with the specific consumers entitled to use the customized coupons.\nIt is an even further object of the invention to provide coupons which are customized based on a certain date, time of day or day of the week, in addition to proximity-based and/or demographic based customization."}
{"text": "During some surgical procedures, most notably hernia repair procedures, it is considered desirable by many practitioners to reinforce the muscle tear or other defect with a piece of surgically implantable mesh. Physicians most often use an open-weave, sintered mesh made of polypropylene and hold it in place by a type of permanent fixation method. One common method of fixation uses metallic fasteners, such as staples, which remain in the body permanently after the hernia repair. Medical device designers have created a number of larger-size devices to fasten tissue and/or surgical materials to tissue during minimally invasive surgery. According to these designs, the devices typically contemplate a 10 mm or wider application tool used to deploy a fastener. For example, Green et al., U.S. Pat. No. 5,356,064, col. 21, 11. 34-53, describes a device for deployment through a 12 mm trocar guide tube wherein the device stacks a set of staples at roughly a forty-five degree angle to the axis of the device to provide greater visibility. See Green et al. FIG. 18.\nHowever, these tools are being judged too large for deployment according to the current minimally invasive techniques which are bringing the size of the surgical instruments and access ports down to a 5 mm diameter. Moreover, Green et al. cannot be readily scaled down because of physical limitations caused by the generally transverse stacking of fasteners. See Green et al., FIG. 18. In addition, the design of Green et al. cannot be effectively scaled down because the fastener discharged by the application tool must be of sufficient scale to securely span across the defect and/or strands of surgical mesh and efficaciously engage sufficient tissue area for adequate gripping strength. Green et al. employs a fastener forming system which unduly reduces the finished span or width of the fastener relative to its initial width. See Green et al., FIG. 20-21, col. 22, 11. 38-48. Thus, reduction in Green et al.'s tool diameter would result in an unsatisfactory gripping area for the finished fastener.\nIn addition to reducing fastener-gripping strength, smaller diameter tools have other problems. For example, miniaturization of the surgical fastening tool increases the likelihood of jamming, a common problem for minimally invasive surgical fastener tools, because the critical tolerances for the device's moving parts would be reduced along with the size of the instrument.\nAccordingly, slight changes in deployment stress and temperature can effect the mobility of the moving tool parts. The Origin Tacker, though of 5 mm diameter, employs a rotational actuation mechanism to deploy a helical fastener. Rotation increases the complexity needed for the actuation mechanism, and creates greater need to ensure reliable translation of trigger action.\nWhat is needed is a space-efficient surgical fastening tool that minimizes its, outer diameter while maximizing the gripping area and strength of the fastener. The fastening mechanism of the desired surgical fastening tool must not be complicated and should be limited to a few actuated parts to reduce the probability of jamming during minimally invasive surgery. The device should be designed to avoid double firing and incomplete firing. The device should also permit for the easy reloading of additional fasteners during extensive surgical procedures. The prior art devices are inadequate to meet these objectives."}
{"text": "1. Technical Field\nThe present invention relates to a system including a plurality of storage devices, and more particularly, to a data transmission technique for the system.\n2. Related Art\nAn ink container which is a detachable liquid container is generally mounted on an ink jet printing apparatus which is an example of a liquid ejecting apparatus. A storage device is installed in some ink containers. The storage device stores various kinds of information such as the amount of ink remaining in the ink container or the ink colors (JP-A-2002-370383 and JP-A-2004-299405). A control device installed in the printing apparatus communicates with the storage device of the ink container.\nIn the known technique, however, the reliability of communication between the control device installed in the printing apparatus and the storage device of the ink container is hardly taken into consideration. For example, a problem may arise in that the communication between the control device installed in the printing apparatus and the storage device of the ink container may fail due to the contact failure of an electric connection portion between the printing apparatus and the ink container. When the printing apparatus continues an operation in the communication failure state, a problem may arise where a fault occurs in the stored details of the storage device.\nIn general, it is desirable that the data transmission speed is as rapid as possible. For example, the data transmission speed can be made more rapid by increasing the clock frequency of data transmission. However, when the clock frequency is excessively increased, a problem may arise in that consumed current flowing in a data line is excessively increased. In an example according to the related art, little consideration was paid to the consumed current necessary for communication (data transmission) between the control device and the storage device. When the consumed current necessary for communication between the control device and the storage device is reduced, current capacity of a data transmission circuit can be reduced. Therefore, since the configuration of the circuit can be simplified, an advantage can be obtained so problems with heating or noise hardly occur. In an example according to the related art, however, the study regarding this point is not sufficiently conducted.\nThe above-mentioned problems occur not only in the printing apparatus but also in the system including a plurality of storage devices."}
{"text": "A local area network operating under the Gigabit Ethernet standard includes transceiver nodes communicating with each other in a full-duplex mode; that is, transmitting in both directions simultaneously at full bandwidth within a particular frequency band. The transceivers typically have a transformer at their output, with a communication cable electrically connecting the transformers of the two transceivers.\nWithin the transceiver, both a transmitter section and a receiver section are coupled to the transformer and, necessarily, to each other. Without additional circuitry, the signal at the input to the receiver will include both a received signal from the communication cable and a signal from the transmitter section. The received signal, having been attenuated in the cable, will typically be at a lower signal strength than the signal from the transmitter section within the transceiver.\nSignal isolation between the transmitter and receiver sections of a transceiver may be obtained with a costly, transformer-based, hybrid electrical bridge circuit. However, the degree of isolation provided by such bridge circuits is dependent upon how closely the terminating impedance of the bridge matches the impedance of the communication cable. Because the impedance of the communication cable is dynamic, over time more or less of the transmitter signal is present at the input of the receiver to interfere with detection of the received signal.\nAn alternative technique for reducing the signal strength of the signal from the transmitter section involves generating within the transceiver a signal to cancel the transmitter signal. A signal that is identical to the transmitter signal, but having inverse polarity, is mixed with the combination of transmitter and received signal in order to cancel some portion of the transmitter signal. As a result, the signal-to-noise ratio of the received signal is improved.\nHowever, the efficacy of this technique depends upon the degree to which the cancellation signal matches the transmitter signal. Differences between the two signals will result not only in some portion of the transmitter signal not being cancelled out, but may also inject additional noise into the signal at the input of the receiver. Merely matching the transmitter input signal is not enough, the non-linearities of the transmitter circuitry must also be matched in order for the technique to be fully effective.\nOther challenges in implementing a Gigabit Ethernet network arise from the transformer coupling of the transceiver nodes. A transformer presents the transmitter circuitry with a characteristic impedance load. Because transmitters are typically implemented as current source/sink circuits, and the communication standards require that signal levels of a certain voltage be produced in the network, the power-handling requirement of transmitters is directly related to the impedance load presented by the transformers. As the impedance load is reduced, the transmitter must source/sink proportionately more current to produce the required signal voltages.\nSince the communication cable is double terminated (i.e., is transformer coupled at both ends of the cable), the two transformers are in parallel and together present to the transmitter half the impedance of a single transformer. This requires a transmitter design that provides twice the power. Furthermore, where a transmitter is implemented with circuitry that only sinks current, the transformer is typically provided with a center tap connected to a voltage source. Because the transmitter is only driving half the transformer, the load impedance is halved again, requiring another doubling of the power handled by the transmitter circuitry.\nCircuit designs typically used in Ethernet transmitters may not successfully scale up to the power levels required in such networks. The amplifiers and circuit sources used may become non-linear at these increased power levels, resulting in degraded functionality or failure to meet communications standards.\nBefore undertaking the DETAILED DESCRIPTION OF THE INVENTION below, it may be advantageous to set forth definitions of certain words and phrases used throughout this patent document: the terms “include” and “comprise,” as well as derivatives thereof, mean inclusion without limitation; the term “or,” is inclusive, meaning and/or; “each” means every one of at least a subset of the identified items; the phrases “associated with” and “associated therewith,” as well as derivatives thereof, may mean to include, be included within, interconnect with, contain, be contained within, connect to or with, couple to or with, be communicable with, cooperate with, interleave, juxtapose, be proximate to, be bound to or with, have, have a property of, or the like; and the term “controller” means any device, system or part thereof that controls at least one operation, such a device may be implemented in hardware, firmware or software, or some combination of at least two of the same. It should be noted that the functionality associated with any particular controller may be centralized or distributed, whether locally or remotely. Definitions for certain words and phrases are provided throughout this patent document, those of ordinary skill in the art should understand that in many, if not most instances, such definitions apply to prior, as well as future, uses of such defined words and phrases."}
{"text": "1. Field of Invention:\nThis invention relates to short range communication devices, processes and systems. More particularly, the invention relates to short range communication devices, processes and systems utilizing radio frequency identification systems with shorter user identification and session setup time in interactive services, and a reader device with transponder functionality.\n2. Description of Prior Art\nBluetoothR is a short-range wireless communication system operating at 2.4 GHz and designed to establish connection between two or more devices operating within a ten-meter communication range, for nominal output power, but can be extended up to 100 meters for 100 mW-output power. The details of Bluetooth are described in “The Bluetooth Special Interest Group, Specification Of The Bluetooth System, Volumes 1 and 2, Core and Profiles: Version 1.1, 22 Feb., 2001, which are fully incorporated herein by reference. To discover another Bluetooth terminal, an inquiry message is transmitted searching for other devices in the vicinity. Any other Bluetooth terminal that is listening by means of conducting an inquiry scan, will recognize the inquiry message and respond. The inquiry response is a message packet containing the responding devices Bluetooth terminal Address (BD_ADDR). The Bluetooth terminal address is a unique, 48-bit IEEE address, which is electronically engraved into each Bluetooth terminal. The address is virtually guaranteed to be completely unique, so much so that it can be reliably associated with the device's user, much as can the user's passport number or social security number. One of the problems of Bluetooth terminals is the long time required for discovery of a Bluetooth terminal, which can be up to 10 seconds. In many applications, the device discovery time is too long and restricts the usage of Bluetooth, particularly for interactive services.\nRadio Frequency Identification (RF-ID) is a wireless system that automatically identifies tracks and manages objects via a fast connection between the object and a RF-ID reader. RF-ID principles are described in a publication entitled “Radio Frequency Identification RF-ID: A Basic Primer”, published by the Automatic Identification Manufacturers (AIM) web site (http: //www.aimglobal.org), Oct. 23, 2001 and fully incorporated herein by reference. The object includes a transponder, active or passive as a RF-ID tag, which when in the presence of an electromagnetic zone created by the reader broadcasts an object identity signal. The reader senses and decodes the broadcast signal to identify the object. The object identity is achieved by a connectionless communication that is a connection without a logical connection between the reader and the object.\nThe rising use of RF-ID tags promotes standardization efforts to establish a broadly applicable communication framework between RF-ID tag readers and RF-ID tags. Such a standard is established by European Computer Manufacturers Association (ECMA) International, Geneva, Switzerland in a near field communication standard (ECMA-340) dated December 2002. The standard defines communication modes for a near field communication interface and protocol (NFCIP-1) to be used for RF-ID tag readers and RF-ID tags. It also defines both passive and active communication modes to realize a communication network using near field communication enabled devices for networked products and also for consumer equipment. The standard specifies, in particular, modulation schemes, codings, transfer speeds, and frame format of a Radio Frequency/High Frequency (RF/HF) interface, as well as initialization schemes and conditions required for data collision control during initialization. Furthermore, the ECMA Standard defines a transport protocol including protocol activation and data exchange methods.\nThe aforementioned RF-ID tag reader and RF-ID tag communication is covered in a first approach, as shown in FIG. 5, by the passive communication mode between a phone 1 and a RF-ID tag 700. Correspondingly in a second approach, the active communication mode relates to communications between RF-ID tag readers. The state of the art communication between RF-ID tags readers will be enlightened in view of FIG. 5., wherein a second portable terminal 2 is provided which is connected to a RF-ID tag reader. The RF-ID tag reader may be realized as a stand-alone unit which may be detachably connected to the portable terminal 2 or which may be embedded in the portable terminal 2. Both the RF-ID tag readers of portable terminal 1 and the portable terminal 2 support the active communication mode which means, that the both RF-ID tag readers are adapted to communicate with each other by employing the aforementioned near field communication interface and protocol (NFCIP-1) standard. During active communication mode one of the participating RF-ID tag readers of the portable terminal 1 and portable terminal 2 actively simulates a RF-ID tag whereas the other one actively simulates a RF-ID tag reader. The data communication is established via the wireless communication link 6, which is a radio and high frequency (RF/HF) communication link, respectively. This wireless communication link 6 provide data flow between the application 400 operated on the portable terminal 1 and an application 410 operated on the portable terminal 410.\nBut in contrast to passive RF-ID tags which are energized by the interrogation signal of the RF-ID tag reader, both participating RF-ID tag readers are always energized independently, i.e. have their own power supply to generate the radio/high frequency (RF/HF) signals used for data communications between those. Therefore, the RF-ID tag reader, which actively simulates a RF-ID tag, may be regarded as an active RF-ID tag.\nDespite the advantages provided by the RF-ID tag readers intercommunication several drawbacks of the aforementioned explanation may be identified. The realization of the RF-ID tag reader—RF-ID tag reader intercommunication requires a great effort in development since the interface and protocol definitions are sophisticated and show a significantly complex functionality which is of course required to ensure interoperability between RF-ID tag readers of different manufacturers. Moreover, the realization of the RF-ID tag reader—RF-ID tag reader intercommunication requires energy supply to each participating RF-ID tag reader, which can be problematic in view of portable devices, powered by batteries or accumulators. And in certain circumstances the operation of a portable device is not wanted or even forbidden for example in airplanes or in hospitals. Assume a RF-ID tag reader embedded in a cellular telephone wherein the embedded RF-ID tag reader supports RF-ID tag reader—RF-ID tag reader intercommunication and the user of this cellular telephone wants to establish intercommunication with a fixed RF-ID tag reader. Due to the fact that the embedded RF-ID tag reader has to be energized, intercommunication is not possible in these aforementioned circumstances since the energizing of the embedded RF-ID tag reader results in an unwanted or forbidden energizing of the cellular phone. Comment: This was a very good point!\nFurthermore, RF-ID tag reader—RF-ID tag reader intercommunication competes with well established short range communication standards such as Bluetooth which are implemented in numerous state of the art portable terminals. The implementation of competing communication solutions offering comparable functionality is not economical, especially when considering that the implementation of support of RF-ID tag reader—RF-ID tag reader intercommunication is a complex undertaking. However, a RF-ID tag reader with transponder functionality, which simulates a RF-ID tag while another RF-ID tag reader simulates a tag reader, provides enhanced communication between terminals, particularly when the terminals are mobile deices.\nWhat is needed in the art is a Bluetooth terminal having (1) shortened terminal discovery and user identification time via a connectionless connection to another terminal or access point using RF-ID, the terminal or access point establishing a session enabling the terminal to conduct transactions with an interactive service application, and (2) an advanced reader device for radio frequency (RF) identification which is at least adapted to communicate with a radio frequency identification transponder and which is able to provide radio frequency identification transponder functionality also during unpowered conditions.\nPrior art related to short range communication systems and RF-ID includes:\n1. PCT Publication WO01/39108 A1, published May 31, 2001, discloses a system that includes a mobile communication unit wherein a RF-ID identification module is added between the radio part and an antenna in a mobile telephone that includes a Bluetooth function. The module comprises a mixer for transposing identification messages from identification devices in a 2.45 GHz RF-ID system of the backscatter type to a baseband for further processing in a computer part of the unit. The unit may communicate with a super ordinate system, via a standard call channel or Bluetooth channel. The unit is also cabled to send information to the identification devices for alerting or transmitting data, wherein modulation and encoding can be effected in accordance with the Bluetooth standard enabling the mixer to be included as an integral part of a standard Bluetooth radio.\n2. PCT Publication WO01/45038 A2, published Jun. 21, 2001, discloses a short-range communication system includes an interrogator, which interrogates multiple transponders located on a person and receives information separately from the transponders. The interrogator and/or the transponders may be held in an article of clothing or in personal effects of the person. The system enables the close worn by a person to communicate with a plurality of items and to obtain useful information. The network system may be configured so that any particular transponder may communicate with another transponder or a group of transponders of other locals to activate functions; cause action to otherwise shared data and/or information.\n3. PCT Publication WO01/45319 A1, published Jun. 21, 2001, discloses a first short-range radio link operating within a first link range. The first link is between a stationary unit and a mobile communication device in a wireless network where the stationary unit transmits an interrogation signal to the communication device, via a short-range communication link operating within a second link range, essentially smaller than the first link range. The mobile communication device receives the interrogation signal from the stationary unit and transmits a response signal, including a unique identification number of the mobile communication device to the interrogation signal. The stationary unit receives the response signal and authenticates the identification number. The stationary unit and the mobile communication device establish a connection, via the first short-range radio link.\n4. US Publication 20030114104, published Jun. 19, 2003 discloses a method and system for identifying when a first device is within a physical range of a second device, the second device is to establish substantive communications with. The method includes sensing which particular first device is within a physical range of the second device, and establishing substantive communications with the particular first device.\n5. JP 2001-278495 2001-09-13, Application No. JP0209325 JP, Filed 20020912, entitled Information Processing Apparatus and Method A1 Published 20030327 discloses an information processing apparatus and method capable of easily and rapidly starting radio communication. When a cellular telephone is placed in the vicinity of a personal computer, a strap ID is provided from a non-contact IC tag of a cellular telephone strap to a reader/writer) of the personal computer. The personal computer has a table of correspondence between the strap ID and the Bluetooth device name of the cellular telephone. Upon acquisition of a strap ID, the personal computer references the correspondence table and identifies the cellular telephone as a communication partner of the Bluetooth communication. The personal computer identifies the cellular telephone by referencing the Bluetooth device name among a plurality of Bluetooth devices including the cellular telephone and performs Bluetooth communication with the cellular telephone.\nNone of the prior art discloses or suggests (1) a Bluetooth terminal having a RF-ID identification for establishing a connectionless communication with a service application and conducting transactions with the application on a session basis where the connection time and user identification time are shortened relative to the standard Bluetooth terminal discovery process for conducting sessions, and (2) a reader device included in or coupled to the terminal and adapted to operate in: (A) a reader operation mode and a transponder operation mode, the reader device supporting the near field communication (ECMA-340) standard; (B) a transponder mode or passive mode during periods of time, within which the reader is not energized and (C) operate automatically in a transponder operation mode during periods of time, within which the reader device is not energized."}
{"text": "1. Field of the Invention\nThis invention relates to a direct conversion receiver, and more particularly to a direct conversion receiver suitable for use with a communication system wherein a transmission frequency and a reception frequency are different from each other and transmission and reception are performed simultaneously, and a transceiver (i.e., transmitter-receiver) that includes a direct conversion receiver of the type mentioned.\n2. Description of the Relates Art\nIn recent years, attention is attracted to a direct conversion receiver for use with a communication system of the W-CDMA (Wide Band Code Division Multiple Access) system or the like wherein a radio frequency signal received by an antenna is directly converted into a baseband signal.\nIn the following, a direct conversion receiver will be described with a case wherein it is used in a W-CDMA system taken as an example. However, before the direct conversion receiver is described, the W-CDMA system itself will be described.\nIn the W-CDMA system, transmission and reception are performed simultaneously, and a transmission frequency and a reception frequency are different from each other. The difference between the transmission and reception carrier frequencies is referred as a carrier frequency interval. According to the existing standards of the WCDMA system, a W-CDMA terminal apparatus, i.e., a transceiver, has the following frequency configuration:\nReception carrier frequency fr:                2,110 MHz to 2,170 MHz,        \nTransmission carrier frequency ft:                1,920 MHz to 1,980 MHz,        \nCarrier frequency interval fs (=fr−ft):                190 MHz.        \nFirst, a conventional example of configuration of a radio unit (i.e., transceiver) used as a terminal apparatus in a W-CDMA system will be described. FIG. 1 shows an example of conventional configuration that adopts a single super-heterodyne system as a configuration of a receiver.\nIn the transceiver shown in FIG. 1, duplexer 16 for separating a transmission signal and a reception signal is connected to antenna 1. Low noise amplifier (LNA) 2, band-pass filter 3 for removing an image frequency signal, mixer 4, band-pass filter 5 for intermediate frequency (IF) and variable gain amplifier (VGA) 6 are connected in series in this order to the output of duplexer 16 for a reception signal, and an output of variable gain amplifier 6 is supplied to quadrature demodulator 7. The transceiver includes first local oscillator 17 for generating a first local signal and second local oscillator 18 for generating a second local signal. The first local signal is supplied to mixer 4.\nQuadrature demodulator 7 includes amplifier 8 for amplifying an input signal to quadrature demodulator 7, phase separation circuit 11 for receiving the second local signal as an input thereto and generating an inphase (I) component and a quadrature (Q) component whose phase is displaced by 90 degrees (i.e., π/2) from that of the in-phase component, multiplier 9 for multiplying an output of amplifier 8 and the in-phase component from phase separation circuit 11, and multiplier 10 for multiplying the output of amplifier 8 and the quadrature component from phase separation circuit 11. Here, π is, of course, the ratio of circumference of circle to its diameter. An output of multiplier 9 is supplied as reception baseband I signal 14 to the outside through band-pass filter 12 for baseband. An output of multiplier 10 is supplied as reception baseband Q signal 15 to the outside through band-pass filter 13 for baseband.\nMeanwhile, on the transmission side, quadrature modulator 26, variable gain amplifier 25 for intermediate frequency, band-pass filter 24 for the intermediate frequency, mixer 23, band-pass filter 22 for removing image components and power amplifier (PA) 21 are connected in series in this order. An output of power amplifier 21 is supplied to a terminal for a transmission signal of duplexer 16. Transmission baseband I signal 35 is supplied to quadrature modulator 26 through band-pass filter 33 for a baseband, and also transmission baseband Q signal 36 is supplied to quadrature modulator 26 through band-pass filter 34 for baseband. Quadrature modulator 26 includes frequency divider 31 for dividing the frequency of the second local signal from second local oscillator 18 by 2 to produce a third local signal, phase separation circuit 30 for receiving an output of frequency divider 31 as an input thereto and generating an in-phase (I) component and a quadrature (Q) component whose phase is displaced by 90 degrees (π/2) from that of the in-phase component, multiplier 29 for multiplying an output of band-pass filter 33 and the in-phase component from phase separation circuit 30, multiplier 28 for multiplying an output of band-pass filter 34 and the quadrature component from phase separation circuit 30, and adder 27 for adding outputs of multiplier 28 and multiplier 29 and outputting a result of the addition as an output of quadrature modulator 26. The transmission baseband I signal has a peak value of a component whose phase is the same as that of the third local signal, and the transmission baseband Q signal has a peak value of a component having a quadrature phase to that of the third local signal.\nIn the transceiver, a reception signal is received by antenna 1 and separated by duplexer 16 which removes a transmission signal component from the reception signal. The separated reception signal is amplified by low noise amplifier 2 and inputted to band-pass filter 3 for removing the image component. The image component is a frequency component at a symmetrical position on the frequency axis to the reception signal with respect to the local frequency. The image component must be removed sufficiently by band-pass filter 3 because it otherwise leaks into the frequency band same as that of the signal when the signal is down converted by mixer 4. The reception signal from which the image component has been removed is mixed with the first local signal and down converted by mixer 4 to produce a reception intermediate frequency signal.\nThe first local signal is generated by first local oscillator 17 as described above, and in the example shown, frequency (first local frequency fLO1) of the first local signal is lower substantially by transmission reception carrier frequency interval fs (=190 MHz) than transmission carrier frequency ft. In other words,\nfirst local frequency fLO1 (≈ft−fs):                1,730 MHz to 1,790 MHz.Accordingly, central frequency frm of the down converted intermediate frequency signal is given byfrm=fr−fLO1≈fr−ft+fs=fs+fs=2·fs,and is substantially equal to 380 MHz which is twice the transmission-reception carrier frequency interval.        \nThe intermediate frequency signal is subjected to band-limitation by band-pass filter 5 and then amplified to a level necessary for quadrature demodulation by variable gain amplifier 6. The amplified signal is subjected to quadrature demodulation by quadrature demodulator 7 with the second local signal to produce a set of two reception baseband signals of an in-phase component (reception baseband I signal) and a quadrature component (reception baseband Q signal).\nThe second local signal is generated by second local oscillator 18 as described hereinabove, and in the present example, the frequency (second local frequency fLO2) of the second local signal is substantially equal to twice the transmission-reception carrier frequency interval fs (190 MHz), that is,\nsecond local frequency fLO2≈2·fs=380 MHz, and is substantially equal to central frequency frm of the intermediate frequency signal.\nIn the inside of quadrature demodulator 7, phase separation circuit 11 generates an in-phase component and a quadrature component using the second local signal, and the in-phase component and the quadrature component are multiplied by the intermediate frequency signal by multipliers 9, 10 to generate respective reception baseband signals.\nThe reception baseband signals are subjected to band-limitation by band-pass filters 12, 13 and sent as reception baseband I signal 14, reception baseband Q signal 15, respectively, to a signal processing circuit (not shown) in the following stage so that data decoding of the reception signal is performed by the signal processing circuit.\nWhile the configuration and operation of the reception side are described, the frequency configuration of the reception side is associated closely with the configuration of the transmission side, and therefore, also operation of the transmission side will be described.\nOn the transmission side, a set of transmission baseband I signal 35 and transmission baseband Q signal 36 which are baseband signals produced by processing transmission data by means of a signal processing circuit (not shown) in the preceding stage are inputted and passed though band-pass filters 33, 34 for transmission baseband. Band-pass filters 33, 34 limit the frequency band of transmission baseband I signal 35 and transmission baseband Q signal 36, respectively. The band-limited transmission baseband signals are inputted to quadrature modulator 26 in which the quadrature modulation of those signals is performed.\nQuadrature modulator 26 uses the third local signal produced by dividing the frequency of the second local signal (=380 MHz) by 2 by means of divider 31 therein. Frequency fLO3 of the third local frequency isfLO3=fLO2/2≈2·fs/2=fs (=190 MHz).In the inside of quadrature modulator 26, phase separation circuit 30 uses the third local signal to generate an in-phase component and a quadrature component. The in-phase component and the quadrature component are multiplied by the transmission baseband I signal and the transmission baseband Q signal by multipliers 29, 28, respectively, and resulting signals are added by adder 27 to generate a transmission intermediate frequency signal. Central frequency ftm of the transmission intermediate frequency signal is substantially equal to fs, that is,ftm≈fs=190 MHz.The transmission intermediate frequency signal is amplified to a necessary level by variable gain amplifier 25 and then, after unnecessary waves outside the transmission band are removed from the transmission intermediate frequency by band-pass filter 24, it is supplied to mixer 23. Mixer 23 mixes the first local frequency and the transmission intermediate frequency signal to effect up conversion of the transmission intermediate frequency signal up to the transmission frequency band. First local frequency fLO1 is originally set to fLO1≈ft−fs, where ft is the transmission carrier frequency and fs is the transmission-reception carrier frequency interval, and apparently a correct transmission signal can be obtained if first local frequency fLO1 is added to central frequency ftm≈fs of the transmission intermediate frequency to effect frequency conversion.\nThe transmission signal produced by the up conversion by mixer 23 is supplied to band-pass filter 22, by which unnecessary waves outside the transmission frequency band such as an image frequency component generated inadvertently by mixer 23 are removed from the transmission signal. Then, the transmission signal from band-pass filter 22 is amplified up to a predetermined transmission output level by power amplifier 21 and transmitted through duplexer 16 and antenna 1.\nIt is to be noted that, where first local frequency fLO1 and second local frequency fLO2 are set in such a manner as described above, only two local oscillators 17, 18 can be used to generate all local signals necessary for transmission and reception. While the present specification uses the term “substantially” and the symbol “≈” like a case wherein first local frequency fLO1 is set so as to be “substantially” equal to ft−fs, this is because, as obvious to those skilled in the art, it is not necessary to set fLO1 accurately to ft−fs and, in order to adjust the frequency to a prescribed transmission frequency band or reception frequency band, strictly the frequency band of the baseband signal itself must be taken into consideration. Therefore, in the present specification, displacement of the frequency is permitted as far as modulation, demodulation and frequency conversion can be performed in accordance with the scheme recited in the present specification.\nThe conventional configuration which uses the single super-heterodyne system is such as described above. Although the conventional configuration operates sufficiently, it has the following problems if it is intended to proceed with development of an LSI (large scale integration) in order to reduce the cost and the number of parts of a radio unit in the future.\n1) In order for the receiver to remove an image component before the input of mixer 4, a steep image removing filter is required as band-pass filter 3. To this end, it cannot be avoided to use a passive element such as a SAW (surface acoustic wave) filter or a dielectric filter. Therefore, band-pass filter 3 is not suitable for formation of an LSI chip.\n2) Also band-pass filter 5 used in the intermediate frequency stage performs channel selection, and a steep passive element such as a SAW filter or a dielectric filter must be used also for band-pass filter 5. Therefore, band-pass filter 5 is not suitable for formation of an LSI chip.\n3) Variable gain amplifier 6 in the intermediate frequency stage is a high frequency circuit, and therefore, it is difficult to incorporate variable gain amplifier 6 in an LSI chip in order to integrate it with the baseband unit.\nOne of possible countermeasures to overcome the problems described above is adoption of direct conversion for the receiver. An example of it will be described with reference to FIG. 2. FIG. 2 is a block diagram showing a configuration of a transceiver which adopts direct conversion for the reception side.\nThe transceiver shown in FIG. 2 has a modified configuration to the transceiver shown in FIG. 1 in that it does not include mixer 4, band-pass filter 5 for intermediate frequency and variable gain amplifier 6 for intermediate frequency but instead includes variable gain amplifiers 19, 20 provided on the output side of band-pass filters 12, 13 for baseband, respectively, and a reception signal having passed through band-pass filter 3 for removing an image frequency component is inputted as it is to quadrature demodulator 7. Also, frequency fLO1 of the first local signal generated by first local oscillator 17 is different, and not the second local signal but the first local signal is supplied to phase separation circuit 11 of quadrature demodulator 7. The configuration of the transmission side and the configurations of second local oscillator 18, antenna 1, duplexer 16 and low noise amplifier 2 are same as those of the transceiver shown in FIG. 1. However, the frequency of the first local signal on the transmission side is different.\nIn particular, the transceiver shown in FIG. 2 is different from that shown in FIG. 1 in that, after a reception signal passes through moderate band-pass filter 3, it is converted into a reception baseband signal immediately by quadrature demodulator 7. Here, a steep image removing filter is not required. Quadrature demodulator 7 uses the first local signal as a local signal for generation of a reception baseband signal.\nThe first local signal is produced by first local oscillator 17 similarly as in the transceiver shown in FIG. 1. However, in the example shown in FIG. 2, frequency fLO1 of the first local signal is substantially equal to reception carrier frequency fr. In particular,\nfirst local frequency fLO1 (≈fr):                2,110 MHz to 2,170 MHz.In the inside of quadrature demodulator 7, the first local signal is used to generate an in-phase component and a quadrature component by means of phase separation circuit 11, and the in-phase component and the quadrature component are multiplied by the reception signal by multipliers 9, 10, respectively, to generate reception baseband signals. Accordingly, the signals outputted from quadrature demodulator 7 have a peak value of a component of the reception signal which has a phase same as that of the internal local signal and a peak value of another component of the reception signal which has a quadrature phase to that of the internal local signal.        \nThe reception baseband signals are subjected to band-limitation by band-pass filter 12, 13 for baseband and amplified up to levels required for variable gain amplifiers 19, 20 and sent as reception baseband I signal 14 and reception baseband Q signal 15, respectively, to a signal processing circuit (not shown) in the following stage so that decoding of the reception data may be performed by the signal processing circuit.\nOn the transmission side, the configuration itself is similar to that shown in FIG. 1. However, since frequency (first local frequency fLO1) of the first local signal is replaced by reception carrier frequency fr, mixer 23 operates in the following manner. In particular, since first local frequency fLO1 is substantially equal to fr and central frequency ftm of the transmission intermediate frequency signal is substantially equal to transmission-reception carrier frequency interval fs, mixer 23 extracts a frequency of the difference between them. In particular, although fLO1−fs≈fr−fs, since fs originally is fs=fr−ft, also in this instance, it is apparent that a transmission signal produced by up conversion by mixer 23 is a correct transmission signal.\nAlthough the configuration of a conventional direct conversion receiver is described above, it has the following problems. The problems listed arise from the fact that the frequency of the first local signal generated by the first local oscillator is substantially equal to the reception carrier frequency. Accordingly, where direct conversion is involved, these problems cannot be eliminated.\n4) The first local signal may possibly be radiated through duplexer 16, antenna 1. The radiated first local signal has a bad influence on another receiver.\n5) The first local signal may possibly leak into a reception signal. In this instance, an unstable dc offset occurs with the reception baseband signal outputted from quadrature demodulator 7 and causes saturation of a variable gain amplifier or data decoding error.\n6) A reception signal sometimes has a very high intensity at a place immediately below a base station of the other party of communication, and operation of the first local oscillator may possibly be rendered unstable by interference of the intense reception signal.\nAs described hereinabove, an ordinary direct conversion receiver has problems in radiation of unnecessary waves to the outside of the apparatus by leak of a local signal, disturbance to a local oscillator by an intense reception signal from the outside and occurrence of a dc offset at an output of a quadrature demodulator by leak of a local signal into a reception signal."}
{"text": "Technical Field\nThis disclosure generally relates to devices and methods for capturing and processing images from an environment of a user, and using information derived from captured images. More particularly, this disclosure relates to devices and methods for selection of a device action based on a detected person.\nBackground Information\nToday, technological advancements make it possible for wearable devices to automatically capture images and store information that is associated with the captured images. Certain devices have been used to digitally record aspects and personal experiences of one's life in an exercise typically called “lifelogging.” Some individuals log their life so they can retrieve moments from past activities, for example, social events, trips, etc. Lifelogging may also have significant benefits in other fields (e.g., business, fitness and healthcare, and social research). Lifelogging devices, while useful for tracking daily activities, may be improved with capability to enhance one's interaction in his environment with feedback and other advanced functionality based on the analysis of captured image data.\nEven though users can capture images with their smartphones and some smartphone applications can process the captured images, smartphones may not be the best platform for serving as lifelogging apparatuses in view of their size and design. Lifelogging apparatuses should be small and light, so they can be easily worn. Moreover, with improvements in image capture devices, including wearable apparatuses, additional functionality may be provided to assist users in navigating in and around an environment, identifying persons and objects they encounter, and providing feedback to the users about their surroundings and activities. Therefore, there is a need for apparatuses and methods for automatically capturing and processing images to provide useful information to users of the apparatuses, and for systems and methods to process and leverage information gathered by the apparatuses."}
{"text": "The term xe2x80x9cmultidrug resistancexe2x80x9d (MDR) describes the phenomenon whereby certain cancerous tumor cells develop a resistance to broad classes of cytotoxic agents when exposed to an individual cytotoxic agent. In other words, after a certain period of treatment with a cytotoxic agent which initially shows efficacy in controlling the growth of the tumor, the tumor develops a resistance not only to the specific agent to which the tumor was exposed, but also to broad classes of structurally and functionally unrelated agents. It has recently been found that MDR tumor cells over express a particular membrane glycoprotein known as p-glycoprotein (xe2x80x9cpxe2x80x9d for permeability). This p-glycoprotein is a member of the superfamily of ATP-binding cassette (ABC) transporters. It is thought that the exposure of the MDR tumor cells to a cytotoxic agent causes the induction of this p-glycoprotein which mediates a reverse transport system located on the tumor cell membrane that pumps the cytotoxic agent, as well as the other broad classes of cytotoxic agents, out of the tumor cell thus providing mutiple drug resistance for the cell.\nP-glycoprotein is not just found in tumor cells. It is also expressed in a variety of normal, non-cancerous, epithelial and endothelial cells including in such tissues as the adreneal cortex, in the brush border of the proximal renal tubule epithelium, on the lumenal surface of biliary hepatocytes, in pancreatic ductules, and in the mucosa of the small and large intestine. For purposes of describing the present invention, the presence of p-glycoprotein in the small and large intestine is of particular interest.\nWhen substances are ingested, they are mixed with digestive substances secreted by the body and are ultimately combined in a mixture in the lumen of the intestine. The lumen of the intestine is in contact with certain special epithelial cells which form the mucosa of the intestine or the intestinal wall. Nutrients and other substances present in the intestinal lumen passively diffuse into these intestinal epithelial cells and later diffuse into the portal circulation which carries the nutrients via the blood stream on to the liver. Thus, nutrients and other substances are absorbed into the body and become bioavailable for use by other tissues in the body.\nThe intestinal epithelial cells, however, do not just operate as a vehicle for passive diffusion of nutrients and other ingested substances. In addition, there are various active transport mechanisms located in the outer membrane of the epithelial cells which actively transport various nutrients and other substances into the cell. It is now thought that one of the active transport mechanisms present in the intestinal epithelial cells is p-glycoprotein transport mechanism which facilitates the reverse transport of substances, which have diffused or have been transported inside the cell, back into the lumen of the intestine. It has been speculated that the p-glycoprotein present in the intestinal epithelial cells may function as a protective reverse pump which prevents toxic substances which have been ingested and diffused or transported into the epithelial cell from being absorbed into the circulatory system and becoming bioavailable. One of the unfortunate aspects of the function of the p-glycoprotein in the intestinal cell however is that it can also function to prevent bioavailability of substances which are beneficial, such as certain drugs which happen to be substates for the p-glycoprotein reverse transport system.\nIt has now been found that, surprisingly, the antihistamines of the present invention are coincidentally also targeted by the p-glycoprotein reverse transport system in intestinal epiothelial cells and therefore are not fully bioavailable. The present invention successfully provides a method for enhancing the bioavailablilty of these antihistamines.\nThe present invention relates to a method of enhancing the bioavailability of a piperidinoalkanol antihistamine of Formula I \nwherein\nR is hydrogen or C1-C6 alkyl,\nor a pharmaceutically acceptable salt or an individual optical isomer thereof, in a patient which comprises co-administering to said patient an effective antihistaminic amount of said piperidinoalkanol and an effective p-glycoprotein inhibiting amount of a p-glycoprotein inhibitor. The present invention further relates to a method of treating allergic reactions in a patient, which comprises co-administering to said patient an effective antihistaminic amount of antihistamine of Formula I and an effective p-glycoprotein inhibiting amount of a p-glycoprotein inhibitor. The present invention also relates composition comprising an effective antihistaminic amount of a piperidinoalkanol antihistamine of Formula I and an effective p-glycoprotein inhibiting amount of a p-glycoprotein inhibitor.\nThe present invention provides a method of enhancing bioavailability of a piperidinoalkanol antihistamine of Formula I \nwherein\nR is hydrogen or C1-C6 alkyl,\nor a pharmaceutically acceptable salt or an individual optical isomer thereof.\nAs used herein, the term xe2x80x9cC1-C6 alkylxe2x80x9d refers to a saturated hydrocarbyl radical of straight or branched chain configuration of from 1 to 6 carbon atoms. Specifically included within the scope of the term xe2x80x9cC1-C6 alkylxe2x80x9d are the hydrocarbyl radicals methyl, ethyl, propyl, isopropyl, butyl, sec-butyl, isobutyl, tert-butyl, pentyl, hexyl, and the like. One skilled in the art would immediately recognize and appreciate that the compounds of Formula I possess a chiral center and as such exist in stereoisomeric forms. The present invention applies to the racemic mixture of these stereoisomeric forms as well as to the isolated individual stereoisomers. The individual stereoisomers can be isolated from the racemic mixture by separation techniques which are well known and appreciated in the art including chromatographic methods and selective crystallization techniques.\nThe compounds of Formula I may exist in their free form or as pharmaceutically acceptable salts. Pharmaceutically acceptable salts of the compounds of Formula I are those of any suitable inorganic or organic acid. Examples of suitable inorganic acids include hydrochloric, hydrobromic, sulfuric, and phosphoric acids. Examples of suitable organic acids include carboxylic acids, such as, acetic, propionic, glycolic, lactic, pyruvic, malonic, succinic, fumaric, malic, tartaric, citric, cyclamic, ascorbic, maleic, hydroxymaleic, dihydroxymaleic, benzoic, phenylacetic, 4-aminobenzoic, 4-hydroxybenzoic, anthranillic, cinnamic, salicylic, 4-aminosalicylic, 2-phenoxybenzoic, 2-acetoxybenzoic, mandelic acid, and sulfonic acids, such as, methanesulfonic, ethanesulfonic, and xcex2-hydroxyethanesulfonic acid. Non-toxic salts of the compounds of Formula I formed with inorganic or organic bases are also included within the scope of this invention and include, for example, those of alkali earth metals, for example, calcium and magnesium, light metals of group IIIA, for example, aluminum, organic amines, such as, primary, secondary or tertiary amines, for example, cyclohexylamine, ethylamine, pyridine, methylaminoethanol, and piperazine. The salts of compounds of Formula I may be prepared by conventional means as, for example, by treating a compound of Formula I with an appropriate acid or base. The preferred pharmaceutically acceptable salt for compounds of Formula I is the hydrochloric acid salt.\nCompounds of Formula I may be prepared as described in U.S. Pat. No. 4,254,129, which is hereby incorporated by reference in its entirety.\nThe preferred compound of Formula I is the compound (xc2x1)-4-[1-hydroxy-4-[4-(hydroxydiphenylmethyl)-1-piperidinyl]-butyl]-xcex1,xcex1-dimethyl benzeneacetic acid, which is also known as fexofenadine, and its individual stereoisomers. Fexofenadine, as the hydrochloric acid salt, has been recently approved by the United States Food and Drug Administration (FDA) for use as the active ingredient in the antihistamine known as Allegra(trademark). Allegra is indicated for the treatment of seasonal allergic rhinitis with recommended dosing at 60 mg B.I.D.\nThe present invention provides a method of enhancing bioavailability of the compounds of Formula I. The co-administration of an effective antihistaminic amount of a compound of Fromula I along with an effective p-glycoprotein inhibiting amount of a p-glycoprotein inhibitor provides an enhanced bioavailability for the compounds of Formula I. Bioavailability of a drug is defined as the degree to which a drug becomes available to the target tissue after administration and is conveniently measured as the total amount of drug available systemically. Typically, bioavailability is assessed by measuring the drug concentration in the blood at various points of time after administration of the drug and then integrating the values obtained over time to yield the total amount of drug circulating in the blood. This measurement, called the Area Under the Curve (AUC), is a direct measurement of the bioavailability of the drug. Alternatively, bioavailability may be assessed for fexofenadine by measuring total urine output of fexofenadine, since it is known that fexofenadine is not significantly metabolized after oral administration.\nThe present invention provides for an enhancement of the bioavailability of the drug of Formula I by co-administration of a p-glycoprotein inhibitor. By co-administration of a compound of Formula I and a p-glycoprotein inhibitor, the total amount of the compound of Formula I is increased over that which would otherwise circulate in the blood in the absence of the p-glycoprotein inhibitor. Thus, co-administration in accordance with the present invention will cause an increase in the AUC of the compound of Formula I over that seen with administration of the compound of Formula I alone.\nAs used herein, the term xe2x80x9cpatientxe2x80x9d refers to a mammal, such as, for example, a human, mouse, rat, dog, cat, and the like, which is in need of treatment for an allergic reaction. As used herein, the term xe2x80x9callergic reactionxe2x80x9d refers to a histamine-mediated allergic disease, such as, for example, seasonal allergic rhinitis, idiopathic urticaria, and the like. Such diseases are generally distinguished by an allergen triggered release of histamine from storage cells in tissues. The released histamine binds certain H1-histamine receptors which results in the manifestation of the well known allergic symptioms such as sneezing, itching skin, itching eyes, rhinorrhea, etc. An antihistamine, such as the compounds of Formula I, will block manifestation of the allergic symptoms caused by release of histamine by blocking the H1-histamine receptors in various tissues in the body, such as in the skin, lungs or the nasal mucosa. Antihistamines, such as the compounds of Formula I, are thus well known and effective treatment for allergic reactions in patients.\nEnhancement of bioavailability of a compound of Formula I will provide for a more efficient and effective treatment of the patient since, for a given dose, more compound will be available at the tissue sites at which the antihistamine blocks H1-histamine receptors than in the absence of this enhanced bioavailability.\nAdministration of the compound of Formula I refers to oral administration. The compound of Formula I may be administered orally in any convenient dosage form including, for example, capsule, tablet, liquid, suspension, and the like.\nAn effective antihistaminic amount of a compound of Formula I is that amount which is effective in providing an antihistaminic effect in a patient. An effective antihistaminic amount will vary between about 1 mg to about 600 mg of a compound of Formula I as a daily dose depending upon the type of disease to be treated, the degree of severity of the disease, the species of patient to be treated, the dosage regimen, and other factors which are all well within the abilities of one of ordinary skill in the medical arts to evaluate and assess. A preferred amount however will typically be from about 10 mg to about 240 mg, a more preferred amount will typically be from about 20 mg to about 180 mg, and a further preferred amount will typically be from about 40 mg to about 120 mg. The most preferred amount of a compound of Formula I will be 60 mg or 120 mg. The above amounts of a compound of Formula I can be administered from once to multiple times per day. Typically, doses will be administered on a regimen requiring one, two or three doses per day with one and two being the preferred. The more preferred doseage and regimen will be 40 mg twice per day, 60 mg twice per day, 80 mg twice per day, 80 mg once daily, 120 mg once daily, and 180 mg once daily with the most preferred being 60 mg twice per day and 120 mg once daily.\nAs used herein, the term xe2x80x9cp-glycoprotein inhibitorxe2x80x9d refers to organic compounds which inhibit the activity of the p-glycoprotein mediated active transport system present in the gut. This transport system actively transports drugs which have been absorbed from the intestinal lumen and into the gut epithelium back out into the lumen. Inhibition of this p-glycoprotein mediated active transport system will cause less drug to be transported back into the lumen and will thus increase the net drug transport across the gut epithelium and will increase the amount of drug ultimately available in the blood.\nVarious p-glycoprotein inhibitors are well known and appreciated in the art. These include, water soluble vitamin E; polyethylene glycol; poloxamers including Pluronic F-68; Polyethylene oxide; polyoxyethylene castor oil derivatives including Cremophor EL and Cremophor RH 40; Chrysin, (+)-Taxifolin; Naringenin; Diosmin; Quercetin; and the like.\nPolyethylene glycols (PEGs) are liquid and solid polymers of the general formula H(OCH2CH2)nOH, where n is greater than or equal to 4, having various average molecular weights ranging from about 200 to about 20000. PEGs are also known as xcex1-hydro-xcfx89-hydroxypoly-(oxy-1,2-ethanediyl)polyethylene glycols. For example, PEG 200 is a polyethylene glycol wherein the average value of n is 4 and the average molecular weight is from about 190 to about 210. PEG 400 is a polyethylene glycol wherein the average value of n is between 8.2 and 9.1 and the average molecular weight is from about 380 to about 420. Likewise, PEG 600, PEG 1500 and PEG 4000 have average values of n of 12.5-13.9, 29-36 and 68-84, respectively, and average molecular weights of 570-630, 1300-1600 and 3000-3700, respectively, and PEG 1000, PEG 6000 and PEG 8000 have average molecular weights of 950-1050, 5400-6600, and 7000-9000, respectively. Polyethylene glycols of varying average molecular weight of from 200 to 20000 are well known and appreciated in the art of pharmaceutical science and are readily available.\nThe preferred polyethylene glycols for use in the instant invention are polyethylene glycols having an average molecular weight of from about 200 to about 20,000. The more preferred polyethylene glycols have an average molecular weight of from about 200 to about 8000. More specifically, the more preferred polyethylene glycols for use in the present invention are PEG 200, PEG 400, PEG 600, PEG 1000, PEG 1450, PEG 1500, PEG 4000, PEG 4600, and PEG 8000. The most preferred polyethylene glycols for use in the instant invention is PEG 400, PEG 1000, PEG 1450, PEG 4600 and PEG 8000.\nPolysorbate 80 is an oleate ester of sorbitol and its anhydrides copolymerized with approximately 20 moles of ethylene oxide for each mole of sorbitol and sorbitol anhydrides. Polysorbate 80 is made up of sorbitan mono-9-octadecanoate poly(oxy-1,2-ethandiyl) derivatives. Polysorbate 80, also known as Tween 80, is well known and appreciated in the pharmaceutical arts and is readily available.\nWater-soluble vitamin E, also known as d-xcex1-tocopheryl polyethylene glycol 1000 succinate [TPGS], is a water-soluble derivative of natural-source vitamin E. TPGS may be prepared by the esterification of the acid group of crystalline d-xcex1-tocopheryl acid succinate by polyethylene glycol 1000. This product is well known and appreciated in the pharmaceutical arts and is readily available. For example, a water-soluble vitamin E product is available commercially from Eastman Corporation as Vitamin E TPGS.\nNaringenin is the bioflavonoid compound 2,3-dihydro-5,7-dihydroxy-2-(4-hydroxyphenyl) -4H-1-benzopyran-4-one and is also known as 4xe2x80x2,5,7-trihydroxyflavanone. Naringenin is the aglucon of naringen which is a natural product found in the fruit and rind of grapefruit. Naringenin is readily available to the public from commercial sources.\nQuercetin is the bioflavonoid compound 2-(3,4-dihydroxyphenyl)-3,5,7-trihydroxy -4H-1-benzopyran-4-one and is also known as 3,3xe2x80x2,4xe2x80x2,5,7-pentahydroxyflavone. Quercetin is the aglucon of quercitrin, of rutin and of other glycosides. Quercetin is readily available to the public from commercial sources.\nDiosmin is the naturally occurring flavonic glycoside compound 7-[[6-O-6-deoxy-xcex1-L-mannopyranosyl)-xcex2-D-glucopyranosyl]oxy]-5-hydroxy-2-(3-hydroxy-4-methoxyphenyl)-4H-1-benzopyran-4-one. Diosmin can be isolated from various plant sources including citrus fruits. Diosmin is readily available to the public from commercial sources.\nChrysin is the naturally occurring compound 5,7-dihydroxy-2-phenyl-4H-1-benzopyran-4-one which can be isolated from various plant sources. Chrysin is readily available to the public from commercial sources.\nPoloxamers are xcex1-hydro-xcfx89-hydroxypoly(oxyethylene)poly(oxypropylene) poly(oxyethylene) block copolymers. Poloxamers are a series of closely related block copolymers of ethylene oxide and propylene oxide conforming to the general formula HO(C2H4O )a(C3H6O)b(C2H4O)aH. For example, poloxamer 124 is a liquid with xe2x80x9caxe2x80x9d being 12, xe2x80x9cbxe2x80x9d being 20, and having an average molecular weight of from about 2090 to about 2360; poloxamer 188 is a solid with xe2x80x9caxe2x80x9d being 80, xe2x80x9cbxe2x80x9d being 27, and having an average molecular weight of from about 7680 to about 9510; poloxamer 237 is a solid with xe2x80x9caxe2x80x9d being 64, xe2x80x9cbxe2x80x9d being 37, and having an average molecular weight of from about 6840 to about 8830; poloxamer 338 is a solid with xe2x80x9caxe2x80x9d being 141, xe2x80x9cbxe2x80x9d being 44, and having an average molecular weight of from about 12700 to about 17400; and poloxamer 407 is a solid with xe2x80x9caxe2x80x9d being 101, xe2x80x9cbxe2x80x9d being 56, and having an average molecular weight of from about 9840 to about 14600. Poloxamers are well known and appreciated in the pharmaceutical arts and are readily available commercially. For example, Pluronic F-68 is a commercially available poloxamer from BASF Corp. The preferred poloxamers for use in the present invention are those such as poloxamer 188, Pluronic F-68, and the like.\nPolyoxyethylene castor oil derivatives are a series of materials obtained by reacting varying amounts of ethylene oxide with either castor oil or hydrogenated castor oil. These polyoxyethylene castor oil derivatives are well known and appreciated in the pharmaceutical arts and several different types of material are commercially available, including the Cremophors available from BASF Corporation. Polyoxyethylene castor oil derivatives are complex mixtures of various hydrophobic and hydrophilic components. For example, in polyoxyl 35 castor oil (also known as Cremophor EL), the hydrophobic constituents comprise about 83% of the total mixture, the main component being glycerol polyethylene glycol ricinoleate. Other hydrophobic constituents include fatty acid esters of polyethylene glycol along with some unchanged castor oil. The hydrophilic part of polyoxyl 35 castor oil (17%) consists of polyethylene glycols and glyceryl ethoxylates.\nIn polyoxyl 40 hydrogenated castor oil (Cremophor RH 40) approximately 75% of the components of the mixture are hydrophobic. These comprise mainly fatty acid esters of glycerol polyethylene glycol and fatty acid esters of polyethylene glycol. The hydrophilic portion consists of polyethylene glycols and glycerol ethoxylates. The preferred polyoxyethylene castor oil derivatives for use in the present invention are polyoxyl 35 castor oil, such as Cremophor EL, and polyoxyl 40 hydrogenated castor oil, such as Cremophor RH 40. Cremophor EL and Cremophor RH 40 are commercially available from BASF Corporation.\nPolyethylene oxide is a nonionic homopolymer of ethylene oxide conforming to the general formula (OCH2CH2)n in which n represents the average number of oxyethylene groups. Polyethylene oxides are available in various grades which are well known and appreciated by those in the pharmaceutical arts and several different types of material are commercially available. The preferred grade of polyethylene oxide is NF and the like which are commercially available.\n(+)-Taxifolin is (2R-trans)-2-(3,4-dihydroxyphenyl)-2,3-dihydro-3,5,7-trihydroxy-4H-1-benzopyran-4-one. Other common names for (+)-taxifolin are (+)-dihydroquercetin; 3,3xe2x80x2, 4xe2x80x2, 5,7-pentahydroxy-flavanone; diquertin; taxifoliol; and distylin. (+)-Taxifolin is well know and appreciated in the art of pharmaceutical arts and is readily available commercially.\nThe preferred p-glycoprotein inhibitor for use in the present invention are water soluble vitamin E, such as vitamin E TPGS, and the polyethylene glycols. Of the polyethylene glycols, the most preferred p-glycoprotein inhibitors are PEG 400, PEG 1000, PEG 1450, PEG 4600 and PEG 8000.\nAdministration of a p-glycoprotein inhibitor may be by any route by which the p-glycoprotein inhibitor will be bioavailable in effective amounts including oral and parenteral routes. Although oral administration is preferred, the p-glycoprotein inhibitors may also be administered intravenously, topically, subcutaneously, intranasally, rectally, intramuscularly, or by other parenteral routes. When administered orally, the p-glycoprotein inhibitor may be administered in any convenient dosage form including, for example, capsule, tablet, liquid, suspension, and the like.\nAn effective p-glycoprotein inhibiting amount of a p-glycoprotein inhibitor is that amount which is effective in providing inhibition of the activity of the p-glycoprotein mediated active transport system present in the gut. An effective p-glycoprotein inhibiting amount will vary between about 5 mg to about 1000 mg of p-glycoprotein inhibitor as a daily dose depending upon the particular p-glycoprotein inhibitor selected, the species of patient to be treated, the dosage regimen, and other factors which are all well within the abilities of one of ordinary skill in the medical arts to evaluate and assess. A preferred amount however will typically be from about 50 mg to about 500 mg, and a more preferred amount will typically be from about 100 mg to about 500 mg. The above amounts of a p-glycoprotein inhibitor can be administered from once to multiple times per day. Typically for oral dosing, doses will be administered on a regimen requiring one, two or three doses per day with one and two being the preferred.\nWhere water soluble vitamin E or a polyethylene glycol is selected as the p-glycoprotein inhibitor, a preferred amount will typically be from about 5 mg to about 1000 mg, a more preferred amount will typically be from about 50 mg to about 500 mg, and a further preferred amount will typically be from about 100 mg to about 500 mg. The most preferred amount of water soluble vitamin E or a polyethylene glycol will be from about 200 mg to about 500 mg. The above amounts of water soluble vitamin E or polyethylene glycol can be administered from once to multiple times per day. Typically, doses will be administered on a regimen requiring one, two or three doses per day with one and two being preferred.\nAs used herein, the term xe2x80x9cco-administrationxe2x80x9d refers to administration to a patient of both a compound of Formula I and a p-glycoprotein inhibitor so that the pharmacologic effect of the p-glycoprotein inhibitor in inhibiting p-glycoprotein mediated transport in the gut is manifest at the time at which the compound of Formula I is being absorbed from the gut. Of course, the compound of Formula I and the p-glycoprotein inhibitor may be administered at different times or concurrently. For example, the p-glycoprotein inhibitor may be administered to the patient at a time prior to administration of the compound of Formula I so as to pre-treat the patient in preparation for dosing with the compound of Formula I. Furthermore, it may be convenient for a patient to be pre-treated with the p-glycoprotein inhibitor so as to achieve steady state levels of p-glycoprotein inhibitor prior to administration of the first dose of the compound of Formula I. It is also contemplated that the compound of Formula I and the p-glycoprotein inhibitor may be administered essentially concurrently either in separate dosage forms or in the same oral dosage form.\nThe present invention further contemplates that the compound of Formula I and the p-glycoprotein inhibitor may be administered in separate dosage forms or in the same combination oral dosage form. Co-administration of the compound of Formula I and the p-glycoprotein inhibitor may conveniently be accomplished by oral administration of a combination dosage form containing both the compound of Formula I and the p-glycoprotein inhibitor.\nThus, an additional embodiment of the present invention is a combination pharmaceutical composition for oral administration comprising an effective antihistaminic amount of a compound of Formula I (the antihistamine) and an effective p-glycoprotein inhibiting amount of a p-glycoprotein inhibitor (the inhibitor). This combination oral dosage form may provide for immediate release of both the compound of Formula I and the p-glycoprotein inhibitor or may provide for sustained release of one or both of the compound of Formula I and the p-glycoprotein inhibitor. One skilled in the art would readily be able to determine the appropriate properties of the combination dosage form so as to achieve the desired effect of co-administration of the compound of Formula I and the p-glycoprotein inhibitor.\nThe antihistamine and the inhibitor may be administered alone or in the form of a pharmaceutical composition in admixture or otherwise in association with one or more pharmaceutically acceptable carriers or excipients, the proportion and nature of which are determined by the solubility and chemical properties of the antihistamine and inhibitior selected, the dosage regimen desired and standard pharmaceutical practice. The antihistamines, while effective themselves, may be formulated and administered in the form of their pharmaceutically acceptable acid addition salts, such as the hydrochloride, for purposes of stability, convenience of cystallization, increased solubility and the like. One form of the pharmaceutical composition according to the present invention is a combination pharmaceutical composition where both the antihistamine and the inhibitor are present in the same dosage form.\nThe pharmaceutical composition may be prepared in a manner well known and appreciated in the pharmaceutical art. The carrier or excipient is pharmacologically inert and may be a solid, semi-solid, or liquid material which can serve as a vehicle or medium for the antihistamine and the inhibitor. Suitable carriers and excipients are well known in the art. The pharmaceutical compositon may be adapted for oral administration in the form of a tablet, capsule, liquid, syrup, wafer, chewing gum, suspension, or the like. These preparations may contain at least 4% of active ingredient, i.e., the percent by weight of the antihistamine and the inhibitor, but may conveniently be varied depending upon the particular form so that the active ingredients make up from about 4% to about 70% of the weight of the unit dosage form.\nTablets, pills, capsules, and the like may contain one or more of the following carriers or excipients: binders such as microcrystalline cellulose, gum tragacanth or gelatin; excipients such as starch or lactose; surfactants such as polysorbate 80, and the like; disintegrating agents such as alginic acid, Primogel(trademark), corn starch, sodium bicarbonate, calcium bicarbonate and the like; lubricants such as magnesium stearate or Sterotex(trademark); glidants such as colloidal silicon dioxide; sweetening agents such as sucrose or saccharin; flavoring agent such as peppermint, methyl salicylate or orange flavoring. Capsules may contain, in addition to the ingredients listed above for tablets, a liquid carrier such as polyethylene glycol or a fatty oil. Tablets and capsules may contain other various carriers and excipients which modify the physical form of the dosage unit, for example, as coatings. Thus, tablets may be coated with sugar, shellac, or other enteric coating agents. A syrup may contain, in addition to the active ingredients, sterile water, sucrose as a sweetening agent, preservatives, dyes, and colorings and flavors. Materials used in preparing these various compositions should be pharmaceutically pure and non-toxic in the amounts used.\nFor purposes of parenteral administration, the inhibitor may be incorporated into a solution or suspension. These preparations should contain at least 0.1% of the active ingredient but may be varied from about 0.1% to about 50% by weight thereof. The amount of the inhibitor should be adjusted in such compositions so that an a suitable dosage will be obtained upon administration.\nThe solutions or suspensions may also include one or more of the following adjuvants: sterile diluents such as water, saline, fixed oils, polyethylene glycols, glycerine, propylene glycols or other synthetic solvents; antibacterial agents such as benzyl alcohol or methyl paraben; antioxidants such as ascorbic acid or sodium bisulfite; chelating agents such as ethylene diaminetetraacetic acid; buffers such as acetates, citrates or phosphates; agents for the adjustment of tonicity such as sodium chloride or dextrose. The parenteral preparations may be enclosed in ampules, disposable syringes or multiple dosage vials made of glass or plastic.\nMore particularly, the combination pharmaceutical composition may be in the form of a tablet, a capsule, a liquid, a suspension, a syrup, and the like. The combination pharmaceutical composition, including in tablet form, may be a simple admixture of the antihistamine, the inhibitor, and any necessary and appropriate carriers and excipients. Alternatively, the composition may be in the form of an admixture of various heterogeneous pellets, beads or other heterogeneous particles which provide an appropriate formulation. In addition, the pharmaceutical composition may be in the form of a multiple compression tablet such as a multilayered tablet or a compression-coated tablet.\nCombination pharmaceutical compositions made up of heterogeneous pellets, beads or particles (hereinafter referred to as xe2x80x9cheterogeneous pelletsxe2x80x9d), or made up of multiple compression tablets, are useful for administration of pharmaceutical compositions which provide for different release characteristics for the antihistamine and inhibitor. For example, these compositions may provide for an immediate release of the inhibitor and a sustained release of the antihistamine, or vice versa. These compositions are prepared according to standard techniques which are well known and appreciated in the art such as those described in U.S. Pat. No. 4,996,061 which is hereby incorporated by reference in its entirety."}
{"text": "1. Field of the Invention\nThis invention is related to electrical connectors and to strain reliefs for securing wires to the connector housing and preventing forces or vibrations on the wires from damaging the terminals or the contact interface between mating terminals.\n2. Description of the Prior Art\nStrain reliefs are commonly employed on electrical connectors to prevent damage to terminals and the contact interface between mating terminals. The simplest form of damage to terminals is the potential damage to a wire termination due to tension applied to the wire. A crimped termination can be damaged in this way and even if the damage is not visually noticeable, the gas tight connection between a terminal and a wire can be adversely affected. Oxides can form and eventually the contact between the terminal and the wire can deteriorate causing an increase in contact resistance and potentially even an open circuit, especially where tension is repeatedly applied to the termination.\nEven where there is no physical damage to the termination between the wire and the contact terminal, vibration transmitted through the wires to the terminals can cause the mating interface between spring loaded mating contact terminals to be unstable or to degrade over time. Vibration can cause the contact point between mating terminals to shift slightly from a point where good metal to metal contact is obtained to a location where oxides have built up. Fretting corrosion can also result. For these and other reasons it is important that a strain relief be provided to isolate external forces from terminals and terminations.\nConventional strain reliefs are located on the rear of electrical connectors. In different applications the same connector may be used in some applications where the wires should exit to the side while in other applications the wires may exit to the rear of the connector. Therefore some conventional electrical connector strain reliefs have provided means for securing wires to the stain relief in at least three orientations. U.S. Pat. No. 4,080,035 is one example of an electrical connector in which cable ties can be used to secure wires to the strain relief.\nU.S. Pat. No. 5,211,706 is another example of a prior art connector in which the wires can exit either to one side of the strain relief or through the rear of the strain relief housing. That strain relief comprises a top and bottom housing with one molded cable tie loop on the bottom half and another on the top half. For applications in which the wire is to exit through the rear a cable tie is inserted around the wires through the loop in the top housing half. For side entry applications, the cable tie is inserted through a loop adjacent the side of the bottom half.\nA common problem with these and other prior art multi-exit strain reliefs is the relative size of the strain relief. In some cases, the strain relief occupies as much space as the connector itself. The strain relief may be wider than the connector or it can be longer than the connector. In each case large strain reliefs cannot be used in applications where space is at a premium, which is generally the case where side exit connectors are used."}
{"text": "Fluid containers store fluid to be supplied to other devices. Fluid containers may include multiple chambers and be removably installed in devices such as image forming apparatuses to supply the fluid thereto. Generally, one or more chambers include regulator units to regulate the flow of the fluid in the fluid container and/or the device.\nThroughout the drawings, identical reference numbers designate similar, but not necessarily identical, elements."}
{"text": "The quality of a rendered image is significantly improved with anti-aliasing. In computer graphics, aliasing is the process by which smooth curves and angled lines appear to be jagged because the resolution of the graphics device or file is not high enough to represent a smooth curve or angled lines. Many conventional anti-aliasing techniques exist to reduce aliasing effects (e.g., jagged lines following boundaries between pixels) in rendering a two-dimensional (2D) image of the three-dimensional (3D) object(s).\nSuper-sampling and multi-sampling are two common conventional anti-aliasing techniques. Super-sample anti-aliasing (SSAA) involves generating multiple samples within a pixel, where each sample is independently computed for coverage and shading. The shaded samples are stored within a frame buffer and blended for display. While super-sampling produces a very accurate and high-quality image, super-sampling is quite expensive because each pixel within a rendered image requires the computational processing of multiple fully shaded samples, and shading is typically the most expensive operation within the graphics rendering engine.\nMulti-sample anti-aliasing (MSAA) techniques are a less expensive technique that uses one fully shaded color value and a coverage mask, rather than multiple fully shaded samples, to generate the multiple samples stored in the frame buffer that are ultimately blended to produce a pixel within a rendered image. This combination of a color and a sample mask is sometimes referred to as a fragment. Multi-sampling is commonly used because of the substantial cost-versus-performance benefit that is typically achieved without a significant loss in overall image quality. Although multi-sampling saves shader processing relative to super-sampling, multi-sampling still requires a frame buffer with multiple samples per pixel and the attendant bandwidth, which can limit application performance.\nA common method to reduce the memory bandwidth required by multi-sample anti-aliasing is referred to as color fragment compression. Instead of replicating the fragment color to each sample marked in the fragment's sample mask, fragment compression stores the fragment colors and their sample masks. As a result, if a color covers N samples, it only needs to be accessed once (along with the coverage information) instead of N times. The coverage information is typically encoded as a fragment number per sample and is usually stored in a separate array called a control plane. If there are four samples per pixel, the control plane stores four 2-bit codes per pixel.\nWith the goal of improving on MSAA (esp. regarding anti-aliasing and bandwidth/storage performance), NVIDIA Corporation introduced CSAA (coverage sampling anti-aliasing), and AMD introduced EQAA (enhanced quality anti-aliasing). Each of these extends MSAA. CSAA and EQAA augment the normal samples with additional coverage samples, which do not store color, depth or stencil values, but rather a link to one of the normal samples. The number of fragment colors that can be represented is limited to the number of normal samples, but anti-aliased coverage is computed based on the total number of normal plus coverage samples. This results in quality similar to MSAA with the total number of samples, provided that not too many triangles or other primitives overlap the pixel.\nThe Detailed Description references the accompanying figures. In the figures, the left-most digit(s) of a reference number identifies the figure in which the reference number first appears. The same numbers are used throughout the drawings to reference like features and components."}
{"text": "1. Field of the Invention\nThe present invention relates to an LPI fuel supply system. More particularly, the present invention relates to an LPI fuel supply system that securely supplies an engine with fuel from a fuel tank.\n2. Description of Related Art\nGenerally, a liquefied petroleum injection (LPI) engine uses an injector to inject liquefied fuel into a cylinder.\nThe operation principle of the LPI engine is schematically explained as follows: an electronic control unit (ECU) receives input signals from sensors to detect the condition of the engine and controls a fuel pump, the injector, and an ignition coil so as to optimize air/fuel ratio and engine performance.\nThat is, the fuel pump is controlled according to the fuel amount that is necessary to supply liquefied fuel to the engine, and the injector sequentially injects the fuel into the cylinders to realize an optimized air/fuel ratio.\nThe LPI system that is provided to the LPI engine includes a fuel pump that is disposed inside or outside the fuel tank so as to supply the liquefied fuel to the engine.\nIf the fuel pump supplies the fuel with uniform flux, the injector injects the liquid fuel. Here, the pressure regulator is disposed so as to uniformly sustain the pressure that is applied to the injector. If the pressure that is formed in the supply line is higher than a control pressure, the passage of the pressure regulator is opened to return the fuel to the fuel tank.\nHowever, when the fuel pump is deteriorated or foreign substances are accumulated in the main filter of the fuel pump, there is a problem that the supply capacity of the pump becomes weakened.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
{"text": "The present invention relates to a heat insulating structure comprising a vacuum insulator that is fixed to a surface and a production process thereof.\nConventionally, inorganic materials such as glass fiber and organic materials such as expanded polyurethane have been used for hot insulators and cold insulators. Glass fiber has excellent heat resistance, but it has a relative thermal conductivity as high as 0.03 through 0.05 kcal/m.h..degree. C. and poor thermal insulation performance. Expanded polyurethane has a thermal conductivity of 0.015 kcal/m.h..degree. C. or so; however, the thermal conductivity is still high in the case of using the material in an insulation box of a freezer that has a temperature inside the box extremely low, for example -90.degree. C. or lower. In such a case the insulation wall would have to be extremely thick to obtain the desired insulation performance.\nThus, recently, vacuum insulators have been used such as disclosed in JP-B 61-17263 (B32B5/18), JP-B 63-35911 (F25D23/06), and JP-B 2-54479 (F16L59/06).\nThese vacuum insulators are produced by first sealing an insulating material that comprises fine powders of silica or perlite or open cell expanded polyurethane in a bag that has a multilayer laminate structure comprising gas barrier films, and then exhausting the gas (air) inside the bag to create vacuum conditions in the bag. Such vacuum insulators achieve a thermal conductivity as low as 0.005 through 0.010 kcal/m.h..degree. C. These insulators make it possible to reduce the thickness of the insulation wall of a freezer, which curtails the installation area, and enlarges the volume inside the box, as well as reducing power consumption.\nOn the other hand, conventional ways of fixing such vacuum insulators to the wall of freezers has been used to make a structure as shown in FIG. 14. In this Figure, reference numeral 1 is the vacuum insulator, mentioned above, that comprises two gas barrier films 2, each of which has laminated internal layers of thermally adherable polyethylene or polypropylene, an aluminum and surface protective layer (as disclosed in JP-B 2-54479), and insulation material 5 inserted between the two gas barrier films 2. The insulation material 5 is made of open cell expanded polyurethane for example, and shown in FIG. 14, by partly removing the gas barrier film 2. The inside of vacuum insulator 1 is evacuated to a predetermined pressure; thereafter, the peripheral part 2A of gas barrier films 2 is heated to make the thermally adherable layers seal together. On the surface of vacuum insulator 1 (the upper outside of the surface protective layer), an adhesive sheet 100 is pasted; or an adhesive material, also shown by reference numeral 100, is precoated on the surface of outer door panel 4 to be insulated (for example, the outer door panel of a freezer), and then vacuum insulator 1 is pasted on the outer door panel 4 sheet by sheet.\nIn such a conventional process, preliminary pasting of an adhesive sheet to vacuum insulator 1 or to the surface to be insulated, such as outer door panel 4 is carried out; or preliminary coating of an adhesive is effected, then the vacuum insulator is pasted onto the surface to be insulated sheet by sheet. Once fixed, adjustment of the postion of the insulation is very difficult in these structures. Thus, the fixation of the insulation at a predetermined position is difficult and handling is troublesome. When the position of the insulation is incorrect, peeling of the insulation may be necessary, and the adhesive has to be again applied. The fixation of vacuum insulation has been very complicated by these procedures. In addition, there is a risk of breaking the gas barrier film 2 when it is necessary to remove the insulation and repaste the door panel."}
{"text": "Generally, a display apparatus may display a corresponding image according to a frames provided by an audio and video (AV) source. However, under the circumstance of displaying static images, the display apparatus may still receive the frames transmitted from the audio and video source continuously, and display the same image. In order to reduce the power consumption when the display apparatus receives the frames to display the static images, a frame buffer for storing an entire frame may be disposed in the display apparatus. Thus, under the circumstance of displaying dynamic images, the display apparatus may display a corresponding image according to a frame provided by the audio and video source, and under the circumstance of displaying static images, the display apparatus may store the static frame in the frame buffer, and the display apparatus displays according to the static frame stored in the frame buffer.\nNevertheless, when the audio and video source provides the dynamic frame again, it is possible that the operation timing of the display apparatus may not synchronize with the operation timing of the frame provided by the audio and video source, and therefore, the image displayed by a display panel may be affected when the display apparatus directly provides the dynamic frame to a driving circuit of the display panel. Thus, providing new dynamic frames to the driving circuit of the display panel without affecting the operation of the display panel has become one of the design concerns for this type of display apparatus."}
{"text": "There are many various types of data acquisition systems available for use on personal computers, including an architecture for plug-in printed circuit boards built to the “PCI” standard for IBM-compatible personal computers. Certain PCI cards can be used as digital sampling systems; however, the fastest available such PCI digital sampling systems are not consistent with regard to their sampling rates. For example, when a PCI data acquisition card transfers data to the personal computer, there is a time interval where incoming samples can be lost. If the time interval to transfer data between the PCI data acquisition card and the personal computer is greater than the sampling rate period, then data that needs to be sampled may be lost during this transfer. The result is that the incoming data would not be correctly sampled at the desired rate.\nAnother limitation of the PCI technology is that the data transfer time is limited to the PCI interface parameters, and since the PCI interface is based upon the architecture of the personal computer itself, this transfer time can vary depending upon the chip set of the PC. The end result in using PCI data acquisition cards is that the “raw” data sampling frequency has an upper limit, and furthermore the sampling rate itself is inconsistent.\nIn addition to PCI digital sampling systems, certain other types of digital data receivers have been disclosed in the patent literature, including U.S. Pat. No. 5,539,782, which discloses a digital data receiver that restores the symbol timing of an analog signal received over a digital communications link. The analog signal is transmitted in a form of a complex signal, and the receiver first demodulates and samples the received analog signal, and then samples the demodulated signal by a frequency that is eight times the frequency of a symbol frequency of the transmitted symbols (digital data). This sampled data is presented to a low pass filter, then sampled again to a four to one ratio by the first sampling section, which outputs a frequency that is twice that of the symbol frequency. The low pass filter output is sampled using a first clock that is applied from a symbol timing restoring circuit. The output data of the “first sampling section” is presented to a matched filter, then to an equalizer circuit, and then is sampled again to a two to one ratio at a second sampling section, which outputs symbol data at the symbol frequency. The second sampling section samples the output data of the equalizer in accordance with a second clock applied from the symbol timing restoring circuit. The output data of the matched filter is presented to a timing data generating section. The average value of the positional information of the first sampled signal is determined, and this information is compared to a positive threshold value. The present sampling position can be adjusted to be faster by a certain magnitude if the average value is greater than or equal to the positive threshold value. A second comparison occurs that compares the average value and a negative threshold value, when the average value was less than the positive threshold value, determined earlier. The sampling position can be adjusted by a certain magnitude when the average value is less than or equal to the negative threshold value. If the average value was greater than the negative threshold value and less than the positive threshold value, then the present sampling position is maintained. The first and second sampling means use clocks generated by the symbol timing restoring circuit, and these clocks are used to create the adjusted sampling position.\nU.S. Pat. No. 4,965,884 discloses a method for data alignment that receives parallel digital data signals and an aligned frame pulse. The frame pulse occurs at a regular interval that is a multiple of the data clock period. A set of six clock signals are provided that are each phase-shifted from one another. The multiple clock signals are used to begin sampling the incoming data signals while awaiting a frame pulse. The incoming data is stored in a register or a memory circuit, and once the frame pulse arrives and is discovered, a decision logic circuit analyzes the position of the frame pulse within the shift register storing the data.\nThis decision logic circuit chooses the three clock samples that are most centered on the frame pulse. The data is then sampled using the three selected clocks, and a “majority value” of the multiple samples is used. The sampled data is then aligned with a locally generated frame pulse. The data that has been selected as having the majority value is then output.\nU.S. Pat. No. 5,572,556 discloses a method for controlling the reproduction of a sampling clock signal. Since transmission frequency clock signals and sampling clock signals are almost never of the same frequency and phase, in the prior art these clock signals typically are aligned by running one of the digitized signals through a buffer memory. Alignment is established between the asynchronous clock signals, using a voltage controlled oscillator (VCO) to make the signal frequencies coincide. However, quantization error is contained in the reproduced sampling clock signal, which can cause jitter problems. The method not only determines the difference in frequency between the incoming transmission path clock signal and the sampling clock signal, but can use the frequency difference and phase difference information to detect a “start timing” of a reproduced sampling clock signal. In this manner, not only is the frequency controlled, but also the phase angle. Once the transmission path clock signal is extracted from the incoming data transmission signal, digital data is first stored in a buffer memory, and the “start timing” is detected for reproducing a sampling clock signal, which is then used for generating a phase control signal. A “reception frequency difference data signal” is created that is related to the difference in frequency between the extracted transmission path clock signal and the reproduced “sampling clock signal,” which is then used to generate a frequency control signal. The sampling clock signal is then reproduced in response to both the phase control signal and the frequency control signal. The result is that digital data can be read out from the buffer memory with both controlled frequency and phase.\nU.S. Pat. No. 5,502,750 discloses a jitter attenuator that generates multiple phase clocks to generate a transmit clock having less jitter than the receive clock. The generated transmit clock has the same average frequency as the receive clock. The receive clock is first divided into a series of write clocks for writing data into an elastic buffer, and the transmit clock is divided into a series of read clocks for reading data from the elastic buffer. A phase selector, using a counter, selects one of the multi-phase clocks to be the transmit clock. The phase of the transmit clock that is closest to one of the multiple read clocks is adjusted by the phase comparison of the read and write clocks used with the elastic buffer. An external oscillator that has a constant frequency is used to generate the plurality of multi-phase clocks that are used for the transmit clock output. The phase, rather than the frequency, is adjusted thereby eliminating feedback to an external voltage controlled oscillator, which allows the jitter attenuator of this invention to be integrated on a single silicon substrate.\nU.S. Pat. No. 5,708,686 discloses a method for receiving digital signals having a constant bit rate using cell-structured asynchronous transmission with pauses of different length between individual cells. The digital signals are written into a FIFO memory until the FIFO is half filled, at which time the digital signals begin to be read out, using a readout clock whose frequency is less than the frequency of the receive clock. During this readout, a signal controls the frequency of the readout clock, and this signal is derived from the loading state of the FIFO memory. The control variable generated based on the loading of the FIFO is fed to multiple voltage controlled oscillators, thereby making it possible to cover a large frequency range used for the output data rate. The time interval for querying the loading state of the FIFO memory can be at rather large time intervals, or a cell clock can be obtained from the received digital signals that is used to stipulate the point in time of such queries. This invention can be used with encoded audio and video signals, including MPEG-encoded signals in which the MPEG signals use a constant data rate. While the coder and decoder circuits can be set to different bit rates, during operation the bit rate remains constant.\nU.S. Pat. No. 5,493,589 discloses a sampling circuit that synchronizes a data stream at a high sampling rate. The samples are received at a high sampling rate and are converted into data samples at a low sampling rate which is related to the high sampling rate, in which the low rate samples are synchronized with the high rate samples. The data is stored in memory, and the read address signals are derived by combining the write address signals with a difference address signal that is generated by a circuit that includes a modulo counter. A sample-and-hold circuit stores the count of the modulo counter and supplies it to an allocating circuit that generates the difference address signal. A decoder controlled by the modulo counter supplies a clock signal at the low sampling rate to a second sample-and-hold circuit that is provided with data from the memory. The second sample-and-hold circuit outputs synchronized data at the low sampling rate.\nU.S. Pat. No. 5,828,698 discloses a processing system which allows Cellular Digital Packet Data (CDPD) to be transmitted over existing cellular telephone networks. One objective of this invention is to convert from serial in-phase (I) and quadrature-phase symbols at the IS-54 standard symbol rate to time-aligned IQ pairs at the CDPD standard symbol rate. This is accomplished by using two different sampling rates: one rate being an integer multiple of 24.3 kHz (the transmission rate of digital voice), and the second rate being an integer multiple of 19.2 kHz (the transmission rate of CDPD). Once the data has been converted, the invention also provides a method for ensuring that the digital signal processor reads the correct number of samples from the FIFO buffers. This achieves virtual real time processing independent of burst arrival time.\nU.S. Pat. No. 5,796,795 discloses a data transferring circuit which aligns a clock with data. The invention relates to high speed processing of synchronous data, typically in a switch for a broadband network. Alignment is achieved by delaying, by approximately the same amount, a reference clock and the data itself through the data processing section of the circuit, and by regenerating a master clock from the delayed reference clock.\nU.S. Pat. No. 4,841,551 discloses a high speed data-clock synchronization processor. A local crystal clock undergoes a series of parallel output delays through a clock-delayed phase generator. The phased clock signals are then sent to the clock-delayed phase processor. After processing, the optimum clock phase is selected at the optimum phase selector. An optimum clock phase is selected with each received data message. As a result, the utilizable bandwidth may be increased and data distortion is minimized so that the number of stations connected to a data bus provided with the data clock from the synchronization processor may be substantially increased.\nU.S. Pat. No. 4,941,156 discloses a linear jitter attenuation circuit. The jitter attenuator utilizes a FIFO data register for receiving and temporarily storing data in a plurality of storage locations. The FIFO register stores data using a Write pointer driven by a Write clock and Read pointer driven by a Read clock to read the stored data. A phase locked loop is used to synchronize the Read and Write pointers, so to maintain the FIFO in a half empty/half full state. The phase locked loop determines the phase difference between the Read and Write clocks. The phase locked loop is driven by a separate oscillator.\nThe prior patents listed above do not solve the problem of an actual data acquisition device that purports to operate at a constant, fast data sampling rate, but in reality does not do so without missing or losing samples. These documents do not provide a solution that offers a front end device for sampling data at a consistent, fast sampling rate, while allowing the actual data acquisition device to sample at a slower, varying rates."}
{"text": "During the manufacture and assembly of electronic display devices such as for example cathode ray tube (CRT) assemblies for computer monitors and television sets, precise mechanical, magnetic and electronic adjustments are required to ensure the electronic display devices provide optimum reproduction image quality. These adjustments include, but are not limited to, focus, purity of color, convergence of beams, color uniformity, geometry and luminance.\nIn today's computer environment, high resolution electronic display devices are required and therefore, during manufacture it is necessary to ensure the electronic display devices provide high resolution and image quality. This has made visual inspection of electronic display devices to evaluate resolution and image quality unsatisfactory.\nFocus and convergence errors in electronic display devices, which decreases resolution and image quality, are a result of deviations in the profile of the electron beams. Therefore, during testing it is necessary to measure electron beam profile so that focus and convergence corrective measures can be taken. Unfortunately, precise measurement of convergence and focus of electron beams in electronic display devices is hampered by the shadow mask and either the phosphor stripes or phosphor dots therein.\nmethods to measure electron beam profile in an electronic display device are known. For example, Minolta has developed a CB-150 system to analyze focus in an electronic display device by measuring electron beam profile. The system includes a processing unit having a built in video generator, a measuring head and a personal computer. In operation, an electron beam is generated within the electronic display device and the light emitted by three phosphor lines within the electronic display device is detected by a CCD camera in the measuring head. The electron beam is then shifted to the left by 80 .mu.m and the light emitted by the three phosphor lines is detected by the CCD camera again. Thereafter, the electron beam is shifted to the right by 160 .mu.m and the light emitted by the three phosphor lines is detected by the CCD camera yet again. By shifting the electron beam and measuring the color output of the phosphor lines, the profile of the electron beam, which is hidden by the separation of the phosphor lines, can be approximated.\nAlthough this system allows the electron beam profile to be measured to evaluate focus in an electronic display device, the system is slow inhibiting electron beam profile measurements to be made in real time. Accordingly, an improved method and system of measuring electron beam profile in an electronic display device is desired.\nIt is therefore an object of the present invention to provide a novel method and system for measuring electron beam profile to evaluate focus and/or convergence in an electronic display device."}
{"text": "Presently, germination of plant seeds of flowering plants and/or of plant spores from non-flowering plants, may generally involve wetting the seed and/or spore to trigger germination. Such a process may generally involve a given seed and/or a given spore, a moisture source, and a growth medium. Often the moisture source may be water that may have wetted the growth medium, which may then retain some of the water. The growth medium may comprise: soil, potting soil, gardening soil, peat, moss, gravel, sand, vermiculite, fertilizer, manure, paper products (e.g. paper towels, napkins, tissues, and the like), fabrics, textiles, linens, and the like. Once the seed and/or the spore may be brought into contact with sufficient moisture, germination may initiate and then subsequent growth may be dictated by temperature, moisture content of the growth medium, and presence of light/dark cycles, the severity of such light, and possibility an applicable growing zone if not using artificial light and/or temperature sources.\nHowever, the process of initiating seed and/or spore germination has historically and presently remains relatively benign, boring and uninteresting. Aside from wetting the seed and/or the spore in a growth medium, there is little to see and/or do.\nThere then is need in the art for devices, apparatuses, methods, and/or systems of initiating germination in a way which may be exciting, interesting—providing a dynamic visual event to observe and possibly specific smells to smell. For example, there is a need in the art for a solid container with a specific shape, that may comprise one or more seeds (spores), wherein the solid container with the specific shape may be triggered to dissolve in an aqueous acid base chemical reaction producing bubbles and fizzing to expose the plant seeds (and/or plant spores); wherein the wetted plant seeds (and/or wetted plant spores) may be wetted to initiate germination. And in some embodiments, the bubbles and fizzing may also emit specific smells (odors and/or fragrances). The specific shape may be a shape not traditionally associated with growing plants and thus it housing seeds (spores) which may germinate after the specific shape dissolves, may offer a pleasant surprise to the user.\nIt is to these ends that the present invention has been developed."}
{"text": "1. Field of the Invention\nThe present invention relates to a shift system for outboard motors, which is configured to detect a vessel operator's shift operation and operate a motor-driven shift actuator in response to the detected vessel operator's shift operation, to thereby perform a shift operation for the outboard motor in a motor-driven manner.\n2. Description of the Related Art\nConventionally, there have been proposed shift systems for outboard motors, which include a motor-driven shift actuator, and are configured to detect a vessel operator's shift operation and operate the motor-driven shift actuator in response to the detected vessel operator's operation, for performing a shift operation in a motor-driven manner (see e.g. Japanese Laid-Open Patent Publications (Kokai) No. 2004-245350 and No. 2004-244003). In the shift system disclosed in Japanese Laid-Open Patent Publication (Kokai) No. 2004-245350, a motor-driven actuator is provided with a manually operable emergency gear, for failsafe function in case of failure of the motor-driven actuator fails.\nHowever, in the shift system disclosed in Japanese Laid-Open Patent Publication (Kokai) No. 2004-245350, the motor-driven actuator is connected to a shift rod via a plurality of reduction gears and shift rod-side gears and the like, and hence it is difficult for the shift system be compatible with a conventional outboard motor of a type for which the shift operation is manually performed using a shift cable.\nOn the other hand, the shift system disclosed in Japanese Laid-Open Patent Publication (Kokai) No. 2004-244003 is designed without giving sufficient consideration to the orientation of a motor output shaft or the disposition of a motor-driven shift actuator, for compactness thereof, which results in an increase in the size of the whole shift system. Further, the shift system has a position sensor projecting upward in an upper part thereof, which increases the vertical dimension of the system. Furthermore, a neutral switch is incorporated in the motor-driven shift actuator, which increases the size of the motor-driven shift actuator."}
{"text": "In a popular typical lawn mower design, the rider sits atop a four wheeled machine powered by an internal combustion engine. One or more blades swing about a vertical axis within a cutter deck attached to the machine, to cut grasses as the mower moves across a lawn or other turf.\nThe present invention is described here in terms of its application to a transmission-steerable commercial riding mower, for which it is especially useful. Such mowers are typically driven by independent hydraulic motors connected to the rear drive wheels. The front wheels are freely pivotable casters. Directional control, or steering, is obtained by differently varying the speed of the two rear wheels. It is consequently important that good traction be sustained between the rear wheels and the earth surface, especially when the mower traverses a slope, since gravity tends to pivot the front of the mower downhill. It is also important that any lawn mower maintain an evenness of grass cut. Cutter decks are typically either rigidly fixed to the underside of the mower, or they are suspended so they \"float\" --that is, they hang at a fixed elevation, but move upwardly upon contact with a rise in the soil surface.\nLawn mowers are generally intended for use on lawns and other surfaces which are level or gently undulating. Therefore, in most typical three and four wheel riding mowers heretofore, the wheels are directly fastened to the frame of the mower machine. However, then any vertical force or displacement on a wheel of the mower is transmitted directly to the frame (or chassis) of the machine. The frame of the mower will tilt and rise and fall according to the effects of one of the wheels. So, under certain conditions the chassis vertical motion will undesirably change the elevation of the cutter deck with respect to the lawn being cut. Another undesirable effect of lifting or tilting of the mower chassis is loss of wheel traction which is critically important in transmission steerable machines. And, undue bumpy up and down motion is discomforting to the operator.\nFor mowers intended for use on rougher terrains, designs have been employed which seek to mitigate the adverse effects of undulations. Among such are front and rear wheel axles which tilt about a center pivot. But, in machines with a center pivot axle, when a wheel rises, the center pivot point and the whole end of the frame necessarily also rise, albeit to a lesser degree than when the wheel is fastened directly to a frame. Therefore, there is a continuing need for further improvements in suspension systems in self-propelled lawn mowers."}
{"text": "Electro-hydraulic brake apply systems are known wherein the pressure applied to a vehicle's wheel brakes is controlled by a electronic unit that evaluates several parameters and delivers a control signal to a hydraulic modulator that sets the wheel brake pressure. A key parameter used to determine the appropriate braking pressure at the wheel brake is the driver's command, delivered as an input on the brake pedal. Generally, a pressure sensor monitors wheel brake pressure providing feedback to the control system for use in setting the braking pressure.\nA typical hydraulic modulator includes a pressure generation mechanism and a means of controlling delivery of the generated pressure to the wheel brakes. This may take the form of a pump and proportional hydraulic valve, a pump with a pair of two way valves or a movable piston variable pressure chamber device. The number and arrangement of these elements included in a braking system is determined by the system layout and selected control scheme.\nWith a movable piston variable chamber device, a piston is driven linearly in a bore to vary the size of a pressure chamber. The pressure chamber is connected to a brake line leading to the wheel brake. For an application of braking pressure, the size of the variable chamber is reduced, to take up compliance in the system, and to increase braking pressure. The piston applies an increased force to the contained non-compressible fluid to apply the brake. To decrease braking pressure, the force on the piston is reduced and when appropriate, the size of the variable chamber is increased to release the brake. With this type of operation it is preferable to latch the piston in position for various operating requirements.\nOne known method of latching the piston includes the use of an electro-magnetic brake that when engaged, prevents rotation of a motor that drives the piston through an interposed drivetrain. An electromagnetic brake type latch uses frictional engagement with the motor's rotor to lock the rotor in position and thereby, prevent movement of the piston. To allow movement of the piston, the electromagnetic brake is disengaged. Engagement and disengagement of the latch is controlled by an electronic unit."}
{"text": "Personal hygiene and comfortable cleansing following use of a toilet facility are important considerations. For many people the most effective means for cleaning themselves, or others in their care such as the elderly or small children, is moistened, but not wet, paper or cloth. Various approaches are known in the prior art for moistening, especially, paper, but all are subject to certain disadvantages and drawbacks. Paper, such as toilet paper, may be moistened by holding the paper under running water from, e.g., a bathroom sink faucet, but this approach results in either spotty application of the water or excessive wetting of the paper and is an unsatisfactory solution.\nIn another approach, water may be sprayed onto the paper from a hand held spray bottle or from a device mounted on the wall and containing a water reservoir. A hand held spray bottle is often inconvenient to use and store, and requires both hands for operation; one to hold the paper and one to activate the spray. The mounted spray devices typically require a pumping action to activate the spray, and while it may be possible for a user to operate the devices known in the prior art with one hand, it is cumbersome to do so and it is difficult to avoid overspray. In both cases the water reservoir must be refilled when empty. Water held in a reservoir is subject to growth of bacteria and algae, which can be detrimental to the users.\nPackaged pre-moistened wipes may be used for cleansing, but such wipes are typically not flushable, and are significantly more costly than conventional toilet paper.\nAccordingly, there remains a need for an effective and easy to use device for moistening paper or cloth, and especially toilet paper that does not require refilling of a water reservoir, avoids overspray of the water intended for moistening, may be unobtrusively and conveniently mounted for easy access."}
{"text": "Blow-molded bottles of thermoplastic synthetic resinous material are, for unavoidable reasons, produced with round bottoms and it is, therefore desireable to apply flat-bottomed base cups to the bottles so that they will stand up. One machine for applying flat-bottomed base cups to bottles with round bottoms is disclosed in U.S. Pat. No. 4,132,584.\nOne object of this invention is to provide an improved machine for applying flat-bottomed base cups to round-bottomed bottles.\nMore specifically, it is an object of this invention to provide such a machine with a loader mechanism adapted to transfer a plurality of bottles received from the blow-molding machine directly and simultaneously to an orbiting conveyor on which the base cups are applied to the bottles. The plastic bottles are formed in a well-known manner several at a time in the chambers of a blow-molding machine from which they are deposited on the loader mechanism.\nThe loader mechanism in the specific embodiment about to be described is a pivoted structure adapted to swing about a horizontal axis to transfer several bottles simultaneously to the orbiting conveyor.\nAnother object of the invention is to provide a machine in which the base cups are applied to the bottles automatically in response to the movement of an orbiting conveyor, preferably by means of a cam rail located beneath the conveyor.\nAnother object is to provide a machine in which an orbiting conveyor supports a plurality of support units of novel construction each comprising a bottle holder and a base cup seat beneath the bottle holder.\nAnother object is to provide a machine having novel means for moving a line of base cups to an orbiting conveyor on which the base cups are applied to the bottles and simultaneously for carrying away the bottles after the base cups have been applied."}
{"text": "The invention relates to a vertically movable framework to receive weights for counterbalancing stage curtains and the like and more particularly to the type of counterweight carriage which is guided in its vertical movements.\nIn theaters where curtains and scenery are moved vertically, counterweight arbors or carriages are utilized to counterbalance the weight of the curtain or scenery or whatever may be moved vertically onto and from the stage. Conventional counterweight carriage designs currently used in the stage and theater industry have tops and bottoms of the arbors which are of cast iron manufacture. Frequently the techniques of casting result in inaccuracies in the alignment of the top of the arbor with the back plate which occasion additional time on the job for grinding, realignment and shimming as may be necessary. Further, inaccuracies in the carriage construction may add friction and poor performance in the guidance of the counterweight assembly. It thus has been evident for some time that an improved counterweight carriage assembly is needed which will provide an effective relatively frictionless means for raising and lowering stage curtains and the like and which at the same time can be installed without undue labor cost or a high skill level requirement. It has occurred to the inventor that the existing problems may be solved through a proper design utilizing heavy duty aluminum extrusions adaptable to a variety of conditions encountered in various theaters."}
{"text": "This invention relates generally to graphical user interfaces more particularly to using a graphical user interface with a pointing device that provides an infrared data communication link between a computer and another device.\nInfrared light (IR) transceivers are a part of many computer systems, peripherals, personal digital assistants (PDA\"\"s), and portable electronic devices such as digital cameras, cellular phones, and handheld scanners. These IR transceivers allow wireless, line-of-sight communication to take place between two devices that are each equipped with an IR transceiver. For example, a digital camera may be equipped with an IR transceiver to allow it to download pictures to a computer also equipped with an IR transceiver. Since the communication is wireless, a user does not have to connect any wires or cables to initiate a download. Furthermore, since there is no dedicated cable, other IR equipped devices may also communicate with the computer using the computer\"\"s single IR port. To facilitate the use of IR technology, several key producers and sellers of IR equipped devices and IR transceivers have teamed up to form an industry standards organization. This organization is The Infrared Data Association (IrDA). IrDA is an International Organization headquartered in Walnut Creek, Calif. that creates and promotes interoperable, low cost infrared data interconnection standards that support a walk-up, point-to-point user model. These standards support a broad range of appliances, computing and communications devices.\nMany of the IR transceivers on devices today are located on the main chassis. In the case of a personal computer, this means the IR transceiver is located on the CPU enclosure. However, because of shrinking offices, and decreasing availability of desk space, many users would like to locate the CPU enclosure off of the desktop. This presents a problem with IR equipped computers because the IR transceiver on the CPU enclosure is no longer in a position that is convenient for line-of sight communication. Furthermore, even when the CPU enclosure is placed on the desktop, other peripherals (such a keyboard) or desktop clutter (such as soda cans, or loose papers) may obscure the IR transceiver or they may prevent the other device from being placed in a convenient place that has a clear line-of-sight path to the IR transceiver on the CPU enclosure.\nA solution to many of these problems is described in a related application titled xe2x80x9cCOMPUTER POINTING DEVICE WITH INFRARED COMMUNICATION BRIDGExe2x80x9d Ser. No. 09/544,527, Hewlett-Packard docket number 10001665-1, filed on or about the same day as the present application and assigned to the same assignee as the present application. This application discloses a computer mouse or other moveable computer pointing device with an integrated IR transceiver.\nHowever, the integration of the IR transceiver into the pointing device does not address the problem of an easy to use and intuitive user interface for transferring data between the other device and the host computer via the IR transceiver in the pointing device.\nAn embodiment of a user interface incorporating the invention involves placing the cursor of a graphical user interface over an icon or other area of the screen. When communications are initiated by the user, then the data sent across the infrared interface is placed in that area of the screen, or, if the area of the screen was an icon, an action indicated by the icon is taken."}
{"text": "Pigments generally exhibit vivid color tone and high coloring power, and they are widely used in many fields. Examples of use applications in which pigments are used include paints, printing inks, electrophotographic toners, ink-jet inks, and color filters. Pigments are thus important compounds indispensable in everyday life at the present time. General properties, classification by use, and other aspects of pigments are described in, for example, “Dispersion Stabilization of Pigment and Surface Treatment Technique/Evaluation”, Technical Information Institute Co., Ltd., 2001, pp. 123-224. In particular, examples of pigments that require high performance, and that are of particular importance in practical use, are ink-jet ink pigments, and color filter pigments.\nAs the coloring material for ink-jet ink, dyes have been used, but they have points to be improved as to water resistance and light resistance. To improve them, pigments have come to be used. Images obtained from pigment inks have remarkable advantages of superior light resistance and water resistance compared with images obtained from dye-based inks. However, pigments are not easily formed uniformly and pulverized into fine-particles of a nanometer size (i.e., monodispersion), which can permeate pores in the surface of paper, and hence pigments in the image are poor in close bindness to the paper. Consequently, improvement thereof has been demanded.\nAn exemplary important application of a pigment includes a color filter used for a CCD sensor. In particular, recently, there is a demand for a reduction in thickness of the color filter along with an increase in the number of pixels of a digital camera. An organic pigment is used in the color filter, and a thickness of the color filter substantially depends upon a diameter of a particle of the organic pigment. Therefore, even in this field, there is a demand for producing fine particles that are stable in monodispersion at a nanometer size level under control.\nA method of producing organic pigment fine particles is roughly classified into a brake-down method of producing fine particles from a bulk substance by crushing or the like, and a build-up method of producing fine particles by particle growth in a vapor phase or a liquid phase (see “Experimental Chemical Lecture, 4th Edition” edited by the Chemical Society of Japan (Maruzen Co., Ltd.), vol. 12, pp. 411-488, etc.). In general, the breakdown method (crushing method) is widely used. However, according to this method, it is difficult to obtain an organic substance with a nanometer size level, the productivity is remarkably low, and substances applicable in this method are limited.\nAs a procedure for compensating for the above disadvantage, encapsulating pigment fine particles with resin has been proposed. (Technical Information Association “Preparation and Dispersion/Aggregation Control of nano-fine particles and evaluation thereof” 2003, Ch. 1, Sec. 4, (various synthesis methods of polymer-nano-fine particles and problems); Kengo YASUI, Kozue SUNOUCHI (Dainippon Ink and Chemicals, Inc.) “Increase in Resolution with Aqueous Pigment Ink” DIC TECH REV, 2004, No. 10, pp. 11-18; Hiroshi HARADA, Sadahiro INOUE (Dainippon Ink and Chemicals, Inc.) “Microcapsulated Pigment Jet Ink” DIC Tech Rev, 2003, No. 9, pp. 1-7; Masao TANAKA (Dainippon Ink and Chemicals, Inc.) “Micro/Nano Capsul Technique, Microcapsulation of Pigment, and Application to Ink Jet” Industrial Material, 2004, Vol. 52, No. 6, pp. 42-45, etc.). According to this method, an attempt is made so as to enhance water resistance and light resistance, and to increase a gloss regarding printing. However, the recent demand for the increase in resolution has not been satisfied sufficiently. This is because those methods depend upon the breakdown method, and fine particles cannot be regulated to a desired minute size. Further, this method requires to fine particles by grinding as much as possible, and to adsorb resins to the particles completely, which takes a great amount of time and energy in terms of the productivity.\nOn the other hand, it is conducted that a pigment is dispersed with a dispersant. In this case, merely by allowing the dispersant to adsorb to the surface of the pigment, long-term stability, water resistance, and viscosity stability are insufficient. Further, when such the pigment is used as ink-jet ink, due to the strong shear force at a time of the discharge from a nozzle, the dispersant desorbs from the pigment, deteriorates, and further clogs the nozzle, etc.\nAs to the above-mentioned subjects, a number of methods have been studied. Examples of the methods include chemical reforming of the surface of a pigment by oxidation or the like, and a chemical bond formation or the like of a pigment surface and a dissoluble group by grafting or the like (see JP-A-H 10-316909, “JP-A” means unexamined published Japanese patent application.). However, according to any of those methods, it is difficult to effect a reaction completely. As a result, quality of an obtained dispersion liquid can not be improved. The increase in a content of a dispersant can also be considered so as to ensure the stability. However, there arise the decrease in printing performance and the excessive increase in viscosity.\nFurther, a method of adding a polymerizable surfactant to a pigment obtained by a breakdown method, followed by emulsion polymerization (see JP-A-2005-97517). However, in the emulsion polymerization disclosed herein, in order to allow a monomer completely to adsorb to the surfaces of pigment fine particles, a great amount of monomer is required. This method also depends upon the breakdown method, so that the size of a pigment is large, and it is necessary to make the pigment uniform in a finer size range. Further, according to this method, the energy consumption amount is large, since a crushing step is required, and it is necessary to switch the step facility so as to effect a polymerization reaction, resulting in low productivity.\nContrary to this, recently, a method of producing organic fine particles with a nanometer size by a build-up method has been studied. For example, there is a production method by a micro-chemical process (see European Patent Publication No. 1516896 A1). However, the recent high demand with respect to ink-jet ink etc. is not satisfied, and there is still a demand for further improvement and development."}
{"text": "This invention relates to the free drying in an oven of molded fiber trays used for the display and sale of food and other products and particularly the way to control the characteristic arch warpage of this product.\nThe exact reason for such warpage is not fully known but one speculation is that warpage has occurred from the intersection drying differential between the rib tops and bottoms. This arching may be pronounced in molded fiber products containing latticed bottom panels.\nThe warpage of molded fiber products has been a problem for many years since the generally used molded products for display and sale of food has had this drawback of warping during free drying. This problem is described in U.S. Pat. No. 2,704,493 of Randall, granted Mar. 22, 1955, and entitled \"Molded Pulp Articles and Process of Manufacture\". In that patent an attempt was made to cope with this problem of warpage by a method of molding, free drying, and then attempting to press the article back into shape to remove the warping due to free drying.\nA later patent relating to warpage control in a display tray which refers specifically to food trays which are of interest further discusses the material of the molded pulp food trays and their warpage. This is U.S. Pat. No. 3,929,564 of Reifers, granted Dec. 30, 1975, and entitled \"Method of Molding Free Dried Pulp Display Tray\" which discloses the making of a free dried, molded pulp tray with an indented peripheral margin which is formed during the molding of the tray.\nU.S. Pat. No. 1,900,427 of Clapp, granted Mar. 7, 1933, and entitled \"Manufacture of Fibrous Articles\" also recognized the warpage problem and mentioned that drying may be effected between heated dies under pressure to effect a compacting of the product which was discussed in U.S. Pat. No. 2,704,493 of Randall introduced above, as an expensive process as opposed to the less expensive method involving free drying. Therefore Clapp in U.S. Pat. No. 1,900,427, further stating that articles moulded from ordinary paper pulp such as, for example, ground-wood pulp generally tend to warp badly and lose their shapes, has attempted to overcome warpage by incorporating materials such as asbestos fiber into the paper pulp which material incidentally would now be prohibited for use where food products are involved.\nWarpage of molded food trays of the fibrous material generaly used for the molding of such trays has been a recognized problem for many years but the present invention discloses a method which does not require a change in composition of the fibrous molded material nor the elimination of free drying."}
{"text": "High molecular weight polyesters are commonly produced from low molecular weight polyesters of the same composition by solid state polymerization. The low molecular weight polyesters which are used in such solid state polymerizations can be prepared by conventional melt polymerizations. Solid state polymerization is generally considered advantageous in that the handling of high molecular weight ultra-high viscosity molten polymers is eliminated during the solid state polymerization phase. Thermal degradation during the solid state portion of the polymerization is also essentially avoided.\nIn melt polymerizations, the major part of the reaction is transesterification due to the fact that thermal pyrolysis nullifies much of the esterification of carboxyl end groups. A much larger percentage of the reaction in a solid state polymerization is esterification. Thus, a high molecular weight can be obtained in a solid state polymerization with a lower carboxyl concentration than can be achieved in a melt polymerization.\nThe low molecular weight polyester prepolymers utilized in solid state polymerizations are generally in the form of pellets, chips, or finely divided powder. Such pellets can vary greatly in size: however, as a general rule, the smaller the size of the pellets of polyester prepolymer the faster the solid state polymerization with proceed. Such polyester prepolymers are generally converted from the amorphous to the crystalline state prior to solid state polymerization in order to raise their sticking temperature. This is done in order to keep the pellets or chips of polyester prepolymer from sticking together as a solid mass in the solid state polymerization reactor.\nIn the solid state polymerization of a polyester prepolymer the polymerization is carried out at an elevated temperature which is below the melting point of the polyester resin. Such polymerizations are normally conducted in the presence of a stream of inert gas or under a vacuum. Solid state polymerizations are normally conducted on a commercial basis in the presence of a stream of inert gas since it serves to remove volatile reaction products and helps to heat the polyester.\nHeretofore, the form of the polyester prepolymer has essentially dictated the type of solid state polymerization process which could be employed in order to convert the low molecular weight polyester prepolymer into high molecular weight polyester resin. For example, it has generally been accepted practice to use polyester prepolymer in the form of pellets or chips in vacuum and static bed processes and finely ground powder in fluidized bed processes. The reason for this is that experience has shown that finely ground powders tend to agglomerate in vacuum processes, resulting in slower polymerization rates and a need to regrind the high molecular weight polyester resin produced. Experience has also shown that, in static bed processes finely ground powders will channel or fissure, resulting in uneven polymerization and prolonged polymerization rate. On the other hand, the use of pellets or chips in fluidized bed processes is not economically feasible in view of the velocity and volume of inert gas needed to suspend the pellets or chips and the size of the equipment required to do so.\nPolyester prepolymers which are in the form of finely divided powders solid state polymerize at faster rates than do polyester prepolymers which are in the form of pellets or chips. However, polyester prepolymers which are in powder form are difficult to handle and generally must be polymerized in fluidized bed processes. Additionally, the high molecular weight polyester resins which are made utilizing prepolymers which are in powder form are also in the form of powders which are more difficult to process into articles of manufacture. For these reasons polyester prepolymers in powder form have not been widely utilized in commercial solid state polymerization techniques."}
{"text": "It can be difficult for a user to determine software entities which contribute, directly or indirectly, to a particular window being displayed using a processing system.\nFor example, this can particularly be the case for webpages such as “popup windows”, as the user may be uncertain whether a particular webpage they are visiting is causing popup windows to be displayed, or whether something else, such as a virus, is causing the popup window to be displayed. In another example, malware might cause the malware's own custom windows to popup on a user's computer.\nSimilar problems also exist when targeted content is presented to a user via a webpage. For example, a processing system may have a stored cookie which causes targeted information to be presented via a webpage to the user. If the user requires untargeted information, the user may need to delete the cookie from the processing system. However, it can be difficult for the user to determine which cookie contributed to the targeted information being presented to the user via a webpage.\nTherefore, there is a need for a method, system, computer program product and/or computer readable medium of instructions which provides content to a user relating to one or more entities which contributed to a window being displayed using a processing system.\nAs used herein a “threat” comprises malicious software, also known as “malware” or “pestware”, which comprises software that is included or inserted in a part of a processing system for a harmful purpose. The term threat should be read to comprise possible, potential and actual threats. Types of malware can comprise, but are not limited to, malicious libraries, viruses, worms, Trojans, adware, malicious active content and denial of service attacks. In the case of invasion of privacy for the purposes of fraud or theft of identity, malicious software that passively observes the use of a computer is known as “spyware”.\nA hook (also known as a hook procedure or hook function) generally refers to a function provided by a software application that receives certain data before the normal or intended recipient of the data. A hook function can thus examine or modify certain data before passing on the data. Therefore, a hook function allows a software application to examine data before the data is passed to the intended recipient.\nAn API (“Application Programming Interface”) hook (also known as an API interception), a type of hook, refers to a callback function provided by an application that replaces functionality provided by an operating system's API. An API generally refers to an interface that is defined in terms of a set of functions and procedures, and enables a program to gain access to facilities within an application. An API hook can be inserted between an API call and an API procedure to examine or modify function parameters before passing parameters on to an actual or intended function. An API hook may also choose not to pass on certain types of requests to an actual or intended function.\nA process is at least one of a running software program or other computing operation, or a part of a running software program or other computing operation, that performs a task.\nA hook chain is a list of pointers to special, application-defined callback functions called hook procedures. When a message occurs that is associated with a particular type of hook, the operating system passes the message to each hook procedure referenced in the hook chain, one after the other. The action of a hook procedure can depend on the type of hook involved. For example, the hook procedures for some types of hooks can only monitor messages, others can modify messages or stop their progress through the chain, restricting them from reaching the next hook procedure or a destination window.\nA kernel refers to the core part of an operating system, responsible for resource allocation, low-level hardware interfaces, security, etc.\nAn interrupt is at least one of a signal to a processing system that stops the execution of a running program so that another action can be performed, or a circuit that conveys a signal stopping the execution of a running program.\nA system registry is a database used by modern operating systems, for example Windows™ platforms. The system registry comprises information needed to configure the operating system. The operating system refers to the registry for information ranging from user profiles, to which applications are installed on the machine, to what hardware is installed and which ports are registered.\nA hash function (i.e. Message Digest, eg. MD5) can be used for many purposes, for example to establish whether a file transmitted over a network has been tampered with or contains transmission errors. A hash function uses a mathematical rule which, when applied to a file, generates a hash value, i.e. a number, usually between 128 and 512 bits in length. This number is then transmitted with the file to a recipient who can reapply the mathematical rule to the file and compare the resulting number with the original number.\nAn entity can comprise, but is not limited to, a file, an object, a class, a collection of grouped data, a library, a variable, a process, and/or a device.\nIn a networked information or data communications system, a user has access to one or more terminals which are capable of requesting and/or receiving information or data from local or remote information sources. In such a communications system, a terminal may be a type of processing system, computer or computerized device, personal computer (PC), mobile, cellular or satellite telephone, mobile data terminal, portable computer, Personal Digital Assistant (PDA), pager, thin client, or any other similar type of digital electronic device. The capability of such a terminal to request and/or receive information or data can be provided by software, hardware and/or firmware. A terminal may comprise or be associated with other devices, for example a local data storage device such as a hard disk drive or solid state drive.\nAn information source can comprise a server, or any type of terminal, that may be associated with one or more storage devices that are able to store information or data, for example in one or more databases residing on a storage device. The exchange of information (ie. the request and/or receipt of information or data) between a terminal and an information source, or other terminal(s), is facilitated by a communication means. The communication means can be realized by physical cables, for example a metallic cable such as a telephone line, semi-conducting cables, electromagnetic signals, for example radio-frequency signals or infra-red signals, optical fibre cables, satellite links or any other such medium or combination thereof connected to a network infrastructure.\nThe reference in this specification to any prior publication (or information derived from the prior publication), or to any matter which is known, is not, and should not be taken as an acknowledgment or admission or any form of suggestion that the prior publication (or information derived from the prior publication) or known matter forms part of the common general knowledge in the field of endeavour to which this specification relates."}
{"text": "The means with which user of computing devices can communicate with each other continues to grow. Where individuals may once have been reached at one or two telephone numbers, users may presently be reached via, for example, multiple telephone numbers (e.g., home, work, and mobile), home and work fax numbers, multiple email addresses, personal and professional web pages, and various social networking profiles (e.g., Facebook, Twitter, etc.).\nGiven the amount of contact information individuals currently maintain, a need exists for a new approach to managing, organizing, and presenting the various ways individuals may be contacted."}
{"text": "Media networks may simultaneously broadcast multiple types of programming streams, such as high definition (HD) and standard definition (SD) programming streams. In general, HD programming may have a 16:9 aspect ratio and a resolution of about 1920 vertical lines of resolution by 1080 horizontal lines of resolution, or 1920×1080 (2.1 megapixels per frame) or 1280×720 (0.9 megapixels per frame). Standard definition (SD) is generally used to refer to signal formats that are lower than HD. Media networks may provide various types of content in the available formats based on user demand, content provider request, network requirements. The different formats require different resources, for example, HD content requires more bandwidth than SD bandwidth. In general, media networks increase resource utilization in order to provide higher profile content, such as transmitting HD content using additional bandwidth. However, increasing resource utilization to accommodate content introduces inefficiencies and decreases the profitability of providing multiple types of programming formats. Accordingly, media networks would benefit from methods configured to efficiently allocate bandwidth for content for transmission to viewers using existing resources."}
{"text": "A conventional position detection method typically involves measuring a position of an object using signals transmitted from a satellite, such as a GPS.\nAnother known position detection method involves receiving signals transmitted from a plurality of access points at a terminal, and calculating the position of the terminal from differences in times of reception of the signals.\nMore specifically, a method for detecting the position of a mobile terminal in a cellular phone system has been proposed which includes calculating the differences in the reception times of signals transmitted from the access points to the mobile terminal (e.g. differences in propagation times of the signals from the respective access points to the mobile terminal T1−T2 and T3−T2), and multiplying the propagation time differences by light velocity to calculate differences in propagation distances of the signals from the mobile terminal to the respective access points, thereby determining the position of the terminal, in the following manner (see patent document 1).D1−D2=c(T1−T2)D3−D2=c(T3−T2)\nAnother method for detecting a position of a node in a wireless LAN system has been proposed which includes calculating differences in reception times of a signal (differences between the reception times at the respective access points Ti−Tl), which has been transmitted from the node (terminal) and received at the access points, and multiplying the reception time differences by light velocity to calculate differences between propagation distances of the signal from the node to the respective access points ({[P−Pi]−[P−Pl]}=c (Ti−Tl), i=2, . . . , n), thereby detecting or determining the position of the node, in the following manner (see patent document 2).    [Patent document 1] JP-A No. 181242/1995    [Patent document 2] JP-A No. 101254/2004"}
{"text": "This disclosure is generally directed to imaging members, photoreceptors, photoconductors, and the like. More specifically, the present disclosure is directed to rigid or multilayered flexible, belt imaging members, or devices comprised of an optional supporting medium like a substrate, an optional undercoat or hole blocking layer usually situated between the substrate and the photogenerating layer, and at least one chelating containing charge transport layer, wherein at least one is from 1 to about 5, from 1 to about 3, 2, one, and the like, such as a first charge transport layer and a second charge transport layer, an optional adhesive layer, and an optional overcoating layer, and wherein at least one of the charge transport layers contains in addition to the chelating agent at least one charge transport component, and a polymer or resin binder, and where in embodiments the resin binder selected for the undercoat layer is a known suitable binder including a binder that is substantially insoluble in a number of solvents like methylene chloride, examples of these binders being illustrated in copending U.S. application Ser. No. 11/593,658, filed Nov. 7, 2006 on Photoconductors Containing Halogenated Binders by Jin Wu et al., the disclosure of which is totally incorporated herein by reference. In embodiments, there is disclosed a photoconductor where the charge transport layer contains a chelating agent which for example, passivates the conductive carbon residue that is intrinsically associated with the charge transport component, such as the hole transport molecules or compounds in the charge transport layer thereby resulting in a reduction in the CDS counts. The conductive carbon residue is believed to primarily result from the processes used to prepare charge transport components, such as aryl amine molecules.\nIn embodiments there are disclosed low charge deficient spots (CDS) photoconductors where the charge transport layer is comprised of at least one charge transport, a polymeric binder and a chelating agent. Also, when present the hole blocking layer can contain in embodiments phenol resins, known hole blocking layer polymers as illustrated in U.S. Pat. No. 6,913,863, the disclosure of which is totally incorporated herein by reference, which discloses a hole blocking layer, a photogenerating layer, and a charge transport layer, and wherein the hole blocking layer is comprised of a metal oxide, and a mixture of a phenolic compound and a phenolic resin wherein the phenolic compound contains at least two phenolic groups, or chlorinated polymeric resins as the binder and a hydrolyzed aminosilane as the electroconducting species since it is believed that the CH2Cl2 insoluble binders prevent or minimize the migration of hole transport molecules from an upper charge transport layer into lower layers, and then into the undercoat or ground plane layer. Examples of chlorinated homopolymers include polyvinylidene chloride, chlorinated polyvinyl chloride and chlorinated polyvinylidene chloride. Examples of chlorinated copolymers include copolymers of vinylidene chloride, chlorinated vinyl chloride, and chlorinated vinylidene chloride with vinylidene fluoride, tetrafluoroethylene, trifluorochloroethylene, hexafluoropropylene, and the like.\nA number of advantages are associated with the disclosed photoconductors, such as for example the formation of minimal charge deficient spots (CDS) which result in undesirable printing defects and where the spots can be generated from the charge transport layer or layers; minimization or prevention of the migration of hole transport molecules or components from one charge transport layer to another layer in the photoconductor, such as the photogenerating layer and the charge transport layer, and more specifically, from the top or upper charge transport layer into lower layers of the photoconductor, such as lower charge transport layers and the lower photogenerating layer thereby permitting less undesirable charge deficient spots in the developed image generated. The photoreceptors illustrated herein, in embodiments, have extended lifetimes; possess excellent, and in a number of instances low Vr (residual potential); and allow the substantial prevention of Vr cycle up when appropriate; high sensitivity; low acceptable image ghosting characteristics; and desirable toner cleanability.\nAlso included within the scope of the present disclosure are methods of imaging and printing with the photoconductors illustrated herein. These methods generally involve the formation of an electrostatic latent image on the imaging member, followed by developing the image with a toner composition comprised, for example, of thermoplastic resin, colorant, such as pigment, charge additive, and surface additive, reference U.S. Pat. Nos. 4,560,635; 4,298,697 and 4,338,390, the disclosures of which are totally incorporated herein by reference, subsequently transferring the image to a suitable substrate, and permanently affixing the image thereto. In those environments wherein the photoconductor is to be used in a printing mode, the imaging method involves the same operation with the exception that exposure can be accomplished with a laser device or image bar. More specifically, the flexible photoconductor belts disclosed herein can be selected for the Xerox Corporation iGEN® machines that generate with some versions over 100 copies per minute. Processes of imaging, especially xerographic imaging and printing, including digital, and/or color printing, are thus encompassed by the present disclosure."}
{"text": "There have been many investigations in recent years of the effects that small-geometry surface configurations can have on turbulent boundary layers. Particular attention has been paid to the provision of so-called riblet surfaces in which an array of small longitudinal ribs extend over the turbulent boundary layer region of a surface in the direction of fluid flow over the surface, to reduce momentum transport or drag. Experimental results indicate that net or surface drag reductions of up to about 7% can be achieved.\nIn a paper \"Drag Characteristics of V-Groove and Transverse Curvature Riblets\" (presented by M. J. Walsh at the Symposium on Viscous Drag Reduction, Dallas, Tex., Nov. 7-8, 1979) reference is made to an early investigation (1966) by Liu, Kline, and Johnston in which drag reductions of 3-4% were obtained by reducing the turbulent bursting rate (i.e. the rate of break-up of the low speed longitudinal vertices or \"streaks\" that characteristically form in a turbulent boundary layer flow close to a wall) by employing rectangular fins. Walsh's paper reports the investigation of a number of alternative rib profiles and states he was able to obtain a maximum drag reduction of 7% using V-groove riblets.\nIn the prior art and later in this specification the word streak is often used to refer to longitudinal vertices, although strictly speaking a streak is merely the visible manifestation of vertices in smoke tunnel experiments.\nThe drag reduction may be associated with the ability of the riblets to limit random spanwise movements of the streaks, as has been suggested by a number of sources, e.g. R. E. Falco (AIAA-83-0377, AIAA 21st Aerospace Sciences Meeting, Jan. 10-13 1983, Reno, Nev.). Johansen and Smith (AIAA-85-0547, AIAA-Shear Flow Control Conference, Mar. 12-14, 1985, Boulder, Colo.) have shown that cylindrical riblets of a smaller height than the V-groove riblets with which Walsh obtained his optimum results, may have the effect of anchoring and/or creating sites of low-speed streaks in a limited region above the wall surface, but their experiments also showed drag increases of 3% to 8%. It has been attempted to follow Johansen and Smith by using blade-like riblets to anchor or create low-speed streaks by the use of riblets and then use V-groove riblets between the streak anchors to control bursting, but no improvement in net drag reduction has been reported (\"Progress in Astronautics & Aeronautics\", Vol. 123; and S. P. Wilkinson, J. B. Anders, B. S. Lazos & D. M. Bushnell \"Turbulent Drag Reduction at NASA Langley: Progress & Plans\", Royal Aeronautical Society International Conference on Drag Reduction by Passive Means, Vol. 1, Page 1, 1987: RAeS; and \"Riblets\" by M. J. Walsh in \"Viscous Drag Reduction in Boundary Layers\", Ed. D. M. Bushnell & J. N. Heffner).\nThe results reported from these and other previous investigations all showed that although the effect of riblets may be beneficial, it is rather small for practical purposes, and that has led to the search for alternative geometries and combinations of techniques offering larger benefits.\nThus it has been proposed (U.S. Pat. No. 4,706,910 of M. J. Walsh et al) to reduce skin friction more effectively by providing a riblet surface with superimposed large-eddy break-up devices (LEBU's) as manipulators for the outer region of the boundary layer, but such an arrangement adds considerable complexity and is susceptible to wear and especially accidental damage. In another proposal (\"On the Drag Reduction of the Shark Skin\", Bechert, Hoppe and Reif, AIAA-85-0546) it has been suggested that it may be appropriate to use so-called vortex generator surfaces which are configured so as to produce tangential mixing using local flow conditions over a ridge-pattern structure analogous to the configuration of scales on a shark's skin. It is claimed that this could lead to greater drag reduction, but it would also be inherently complicated and difficult to build up a simulation of the scaly surface that the investigators have concluded is responsible for the low drag characteristic.\nIn contrast to these attempts to reduce skin friction by rather complex means, riblet surfaces are less susceptible to damage and they can be formed relatively straightforwardly, e.g. by machining or pressing or by the application of a preformed, e.g. extruded layer. It would clearly be preferable to employ such surfaces if they were able to give greater reductions of skin friction drag.\nThe mechanisms by which riblets influence the turbulent boundary layer over a wall surface is not yet fully understood and it is perhaps significant that different researchers have published conflicting data about the variation of net reduction of skin friction with rib length. If it can be postulated that they bring a turbulent near-wall flow closer to laminar flow conditions, their efficiency is only a fraction of what it might ideally be, bearing in mind that in equivalent turbulent and laminar flow regions the skin friction drag with laminar flow would be 80% less than that with turbulent flow. If a more significant proportion of this potential improvement could be achieved, whether by riblets alone, or by the use of riblets in combination with other means, the use of riblet surfaces could be considerably enhanced."}
{"text": "1. Technical Field\nThe present disclosure generally relates to a touch screen panel integrated display device.\n2. Description of the Related Art\nA touch screen panel is an input device for enabling a user to input a command. For example, a user can input a command by selecting icons displayed on a screen of an image display device in the touch screen panel, using a hand or an object. In some instances, the touch screen panel may replace a separate input device (such as a keyboard and a mouse) connected to an image display device. By integrating the image display device with the touch screen panel, the range of applications for the touch screen panel can be expanded.\nThe touch screen panel may be implemented using different sensing schemes, for example, a resistive scheme, a light detection scheme, or a capacitance scheme. A capacitive touch screen panel includes detection electrodes arranged in an active touch area, and is configured to detect a change in capacitance between the detection electrodes when a hand or an object contacts the touch screen panel. The change in capacitance is used to calculate a touch position of the hand or the object.\nThe touch screen panel may be separately manufactured and attached onto a surface of a display panel. The touch screen panel may also be formed by directly patterning the display panel. A touch screen panel integrated display device, in which the touch screen panel is combined with the display panel, is widely used.\nHowever, in some instances, the detection electrodes in the touch screen panel integrated display device may extend into a non-display area so as to maintain the capacitance between the detection electrodes at an edge of the display area. However, when a gate driving circuit is positioned in the non-display area of the display panel and overlaps with the detection electrodes, a coupling effect between the detection electrodes and the gate driving circuit may result in noise and generate a ghost touch, thereby degrading touch performance."}
{"text": "The present invention relates to insulated chimney pipes of the type that are commonly used to conduct hot waste gases from a domestic sold fuel burner, oil or gas central heating boiler or the like to the atmosphere in place of traditional brick chimneys.\nSuch chimney pipes must satisfy stringent British Standards tests administered by the Agrement Board to ensure that they satisfy requirements as to ease of installation and access, adequate service life and low thermal conductivity. Chimney pipes frequently pass through floors and lofts and may thus pass very close to combustible material. It is therefore important that the external temperature of chimney pipes can not rise above a predetermined level, and in order to satisfy the relevant British Standard on this criterion most insulated chimney pipes are of a double wall construction. The inner wall is invariably of stainless steel whilst the outer wall may be of galvanised or vinyl coated steel and spaced from the inner wall by a layer, typically about 40 mm thick, of thermally insulating material such as mineral wool, fire-clay or rock granules. However, for chimneys of a larger diameter even a double wall construction may not be adequate. Thus a 200 mm diameter chimney for use with open fires must satisfy British Standard No. 4543, and in order to fulfil the thermal conductivity criterion it has been found necessary to make this chimney of a three wall construction, which is naturally extremely expensive.\nSuch conventional two or three wall chimneys with a stainless steel inner wall are intended to have a service life of about 20 years. In practice this service life is not always achieved because of the highly corrosive atmosphere to which the stainless steel is exposed. In particular the action of water formed by condensation and sulphuric acid formed by the combination of this water with sulphur containing gases liberated by the combustion of sulphur containing fuel oils or the like together with the cumulative corrosive effect of thermal shock or stress can lead to a substantial reduction in the service life and in extreme cases the total collapse of the stainless steel wall. The recent increase in the burning of wood in domestic boilers and open fires has exacerbated these problems since the combustion of wood produces wood acids and other particularly corrosive substances which over time attack and ultimately destroy the stainless steel inner wall of the lining."}
{"text": "1. Field of the Invention\nThis invention relates to locking swivel apparati. In some embodiments, the locking swivel apparati may be used in fishing to fasten terminal gear to a main line.\n2. Description of the Related Art\nIn some embodiments, the locking swivel apparati of the invention may be used for fishing. The practice of fishing typically entails use of a rigid snap connector at the end of a primary line. Such snap connectors are generally devices for attaching fishing lines by a wire snap or clip. These wire snaps are generally small closeable wire structures with an openable loop for attaching gear to the end of a primary line. For example, devices such as double end locking wire snaps, twisted wire hook closures with eyelets, and devices with dual locking arms with eyelets have been used for some time.\nIn many situations a swivel body can also be used in attachment of lines. Such swivel bodies generally comprise an upper end to which a main line can be attached, and a lower end configured to rotate relative to the upper end. Terminal gear can then be connected to the lower end of the swivel by a wire snap in order to allow the terminal gear to rotate without twisting the main line.\nAnglers on larger rivers and lakes frequently fish for various species in varying water depths with a variety of fabricated spoons, lures, bottom bouncers, weights, and natural bait hook arrangements. Many of these devices function correctly with only a non-swiveling snap connector to the line, while others also require a swivel to avoid twisting of the gear or the line and to provide proper action for the hook assembly. Many existing swivel devices also include a snap. However, one must remove the entire assembly in order to replace it with a snap when a non-swiveling snap is desired."}
{"text": "So-called anatomical prostheses of the shoulder are known, which reproduce naturally the gleno-humeral anatomy, and comprise a glenoid insert anchored to the scapula by means of a support attached to the glenoid cavity; the insert articulates with a spherical element attached to the top of the humerus by means of a supporting body and an attachment rod inserted into the humerus itself.\nSo-called inverse prostheses of the shoulder are also known, which reproduce the gleno-humeral anatomy in an inverse manner, and which, unlike anatomical prostheses, comprise a spherical element attached to the scapula and a mating humeral cup or insert attached to the upper end of the humerus.\nIn order to anchor the support to the glenoid cavity it is usually provided to insert into the bone seating a rod substantially passing through the axis of the support, and to use two or more screws that attach to the bone seating a platelet solid with the rod. The screws are normally disposed on the periphery of the rod, and are attached on the scapula substantially along the same direction as the rod.\nDue to particular bone pathologies, or the anatomical morphology of the scapula, or the lack of a sufficient bone thickness, known attachment solutions are not always able to guarantee a correct and enduring attachment of the support to the glenoid cavity. This may entail possible mechanical yielding of the prosthetic part constrained to the glenoid cavity, and consequent problems, obviously, for the patient concerning the progressive loosening and possible reciprocal detachment of the components of the prostheses.\nMoreover, even in substantially normal bone conditions, it may be difficult for the surgeon, given the extremely limited space available, to intervene with medical instruments in order to screw the traditional screws inside the bone of the scapula, and this may cause a non-optimum attachment of the support.\nEP-A1-1.488.764, in the name of the Applicant, discloses an inverse prosthesis able to be disassembled and comprising a double-conical pin element able to be used for mounting a spherical head onto a glenoid insert attachable on the glenoid cavity of the shoulder.\nEP-A1-1.437.099 discloses a device for protecting and covering a femoral neck, which comprises a hollow sleeve able to be mounted on the reshaped femoral neck, and a rigid platelet or flange which extends outwards the sleeve and has a plurality of through holes for the passage of a relative fastener fixable to the bone.\nHowever, these documents do not provide a solution for improving the attachment conditions of a prosthesis to a glenoid seating of a shoulder, more particularly for guaranteeing a correct and enduring attachment.\nOne purpose of the present invention is therefore to achieve a prosthesis of the shoulder, both of the anatomical or inverse type, suitable to guarantee an optimum attachment thereof to the scapula, substantially in any pathological and anatomical condition of the patient, or operating condition of the surgeon.\nAnother purpose of the present invention is to resolve the problems of the state of the art, in a simple and economical manner, without in any way affecting the correct functionality of the prosthesis.\nAnother purpose is to allow to modify, also over time, the conditions of attachment, for example as the consequence of a progressive bone degeneration, without intervening on the normal components of the prosthesis.\nThe Applicant has devised, tested and embodied the present invention to overcome the shortcomings of the state of the art and to obtain these and other purposes and advantages."}
{"text": "There are various types of commercially available orthoses that are intended for stabilizing and immobilizing the lumbar and sacral region of the body for the treatment of low back pain due to lumbar sprain, spinal stenosis, disc herniation, and different degenerative spinal disorders. However, presently available lumbo-sacral orthosis suffer from one or more shortcomings or deficiencies.\nNo two human bodies are the same. Even among individuals with similar weight and body types, a lumbo-sacral orthosis that may serve one individual well, but may be ill fitting on another. Accordingly, there is a need for a lumbo-sacral orthosis that is easily adjustable to ensure a proper fit and to maintain substantial immobilization of the lumbar and sacral region of the wearer at all times, including while standing, sitting or walking.\nLumbo-sacral orthoses are typically prescribed by orthopedic doctors and fitted to patients by a fitting practitioner during an office visit to ensure a proper fit. However, if the patient removes the orthosis for bathing or other reasons, it is difficult for the patient to position and adjust the orthosis to achieve the same fit as when it was initially fitted by the fitting practitioner. Accordingly, there is a need for a lumbo-sacral orthosis which a wearer can replace after removal and achieve the same or substantially the same fit as when initially fitted by fitting practitioner.\nAlso, as a patient moves during the course of the day, commercially available lumbo-sacral orthosis tend to ride-up on the wearer when moving from standing position to a sitting position, particularly if the orthosis includes an abdominal panel. For example, when a wearer sits, the abdominal panel will often contact the wearer's lap causing the orthosis to be forced upwardly causing discomfort to the wearer. Additionally, when the wearer returns to a standing position, the orthosis may no longer be properly positioned and in full contact with the wearer's body. If the orthosis is not in full contact with the wearer's body, the wearer's lumbar and sacral region will not be properly supported or substantially immobilized, permitting the wearer to slouch causing further injury or delaying recovery. Accordingly, there is a need for an lumbo-sacral orthosis that will remain in place as the wearer moves throughout the day, remains comfortable and in full contact with the wearer's lumbar and sacral region and will permit the wearer to sit without causing the orthosis to be forced upwardly and out of proper position."}
{"text": "In general, automatically shiftable transmissions include a number of gear elements and selectively engageable friction elements (referred to herein as torque transmitting mechanisms) that are controlled to establish one of several forward speed ratios between an input shaft and an output shaft of the transmission. The input shaft is coupled to an engine through a fluid coupling such as a torque converter, and the output shaft is coupled to the drive wheels of the vehicle through a differential gearset.\nShifting from a currently established speed ratio to new speed ratio involves, in most cases, disengaging a torque transmitting mechanism (off-going torque transmitting mechanism) associated with the current speed ratio and engaging a torque transmitting mechanism (on-coming torque transmitting mechanism) associated with the new speed ratio. Each such shift includes a preparation or fill phase during which an apply chamber of the on-coming torque transmitting mechanism is filled with pressurized fluid in preparation for torque transmission. Once filled to a predetermined amount, the fluid pressure supplied to the torque transmitting mechanism is lowered to allow smooth torque transmitting mechanism engagement.\nTypically, the fill phase is carried out by commanding a maximum on-coming torque transmitting mechanism pressure for an empirically determined fill time. See, for example, the U.S. Pat. No. 4,707,789 to Downs et al., issued on Nov. 17, 1987, and assigned to the assignee of the present invention and hereby incorporated by reference in its entirety. In Downs et al., the fill time for any given torque transmitting mechanism is obtained from a look-up table as a function of the commanded fill pressure, less the torque transmitting mechanism return spring pressure, and then adjusted as a function of fluid temperature to compensate for variations in fluid viscosity. In operation, the stored fill times are then adaptively adjusted based on the observed inertia phase delay, in order to compensate for sources of error that influence the fill time. A similar approach is found in the U.S. Pat. No. 5,216,606 to Lentz et al., issued on Jun. 1, 1993, and assigned to the assignee of the present invention and hereby incorporated by reference in its entirety. In Lentz et al., the fill time is determined by table look-up based on the pump speed, as compensated for efficiency and fluid viscosity, and further modified based on the time since the last shift and a hydraulic restriction factor.\nTraditional torque transmitting mechanism engagement methods provide a trim phase subsequent to the fill phase to allow smooth engagement and enhance shift quality. The trim phase provides a steady increase in torque transmitting mechanism apply pressure to the on-coming torque transmitting mechanism. Pressure variations or control system instability during the trim phase may cause poor shift quality."}
{"text": "Various communication systems having a plurality of network servers may provide connectivity access to a service for one or more users. The connection may be performed, for instance, using a single device or any number of devices (e.g., mobile phones, personal digital assistants, computers, internet tablets, etc.). The network servers may be located in different terrestrial locations around the globe or be extraterrestrial."}
{"text": "Many people desire to have easy access to their wireless devices at any time—even when exposed to cold environments. Accordingly, there is a need for improved products for allowing users easy access to their wireless devices."}
{"text": "Brevibacillus is a spore-forming bacterium that has been suggested for probiotic effects. For example, Brevibacillus brevis is now well established as a biocontrol agent in many areas, and has been shown to have efficacy against Botrytis and powdery mildew disease (Edwards and Seddon, Edwards and Seddon (1992) Recent Advances in Botrytis Research. The Netherlands: Pudoc Scientific Publications; Bacillus brevis as a Biocontrol Agent against Botrytis cinerea on Protected Chinese Cabbage; pp. 267-271) and Fusarum head blight (FHB) (Zhang et al. (2005) J Zhejiang Univ Sci B. 6(8):770-777). By comparing the activity of B. brevis Nagano against Botrytis cinerea with that of pure gramicidin S and the antibiotic-negative mutant B. brevis E-1, Edwards and Seddon ((2001) J Appl Microbiol. 91:652-659) showed that the mode of antagonism exhibited was antibiosis due to the presence of gramicidin S. There are some other antibiotics (for example tyrocidins and gramicidin D) reported to be produced by B. brevis (Saito et al., 1970, Adv Enzymol. 33:337-380).\nBrevibacillus laterosporus comb. nov. (Shida (1996) Int. J. Syst. Bacteriol. 46:939-946), previously classified as Bacillus laterosporus, is an aerobic spore-forming bacterium that can also demonstrate pathogenicity to insects. In common with B. sphaericus and B. thuringiensis, B. laterosporus produces parasporal bodies, which in this species may be canoe-shaped and which serve to cradle the spore (Hanney (1957) J. Biophys. Biochem. Cytol. 3:1001-1010) or can even be present in different shapes (Smirnova et al. (1996) Res. Microbiol. 147:343-350). However, these parasporal bodies were not considered to have any entomocidal activity (Favret and Yousten (1985) J. Invertebr. Pathol. 45:195-203) until Orlova et al. ((1998) Appl. Environ. Microbiol. 64:2723-2725) demonstrated that some crystals produced during sporulation are highly toxic to Aedes aegypti and Anopheles stephensi larvae.\nSome B. laterosporus strains show no apparent toxic activity to any test organism, and the observed toxicity is not homogeneous among toxic isolates (Favret and Yousten (1985) J. Invertebr. Pathol. 45:195-203, Rivers et al. (1991) J. Invertbr. Pathol. 58:444-447, and Singer (1996) Adv. Appl. Microbiol. 42:219-261). The results of the first bioassays with B. laterosporus demonstrated that some strains presented a larvicidal activity which was 1,000 times lower than that of the B. thuringiensis var. israelensis standard (Favret and Yousten (1985), Rivers et al. (1991)). These results discouraged the use of B. laterosporus in biological control. No entomocidal activity has been demonstrated against plant pathogens using isolates of Brevibacillus sp."}
{"text": "Businesses worldwide recognize the commercial value of their data and seek reliable, cost-effective ways to protect the information stored on their computer networks while minimizing impact on productivity. Protecting information is often part of a routine process that is performed within an organization.\nA company might back up critical computing systems such as databases, file servers, web servers, and so on as part of a daily, weekly, or monthly maintenance schedule. The company may similarly protect computing systems used by each of its employees, such as those used by an accounting department, marketing department, engineering department, and so forth.\nGiven the rapidly expanding volume of data under management, companies also continue to seek innovative techniques for managing data growth, in addition to protecting data. For instance, companies often implement migration techniques for moving data to lower cost storage over time and data reduction techniques for reducing redundant data, pruning lower priority data, etc.\nEnterprises also increasingly view their stored data as a valuable asset. Along these lines, customers are looking for solutions that not only protect and manage, but also leverage their data. For instance, solutions providing data analysis capabilities, improved data presentation and access features, and the like, are in increasing demand."}
{"text": "1. Field of the Invention\nThe present invention relates to a power supply device which charges a main capacitor for use in a flash apparatus or the like.\n2. Description of Related Art\nIn the field of cameras having electronic flash apparatuses, it has become necessary to reduce the sizes of component parts of flash apparatuses as the sizes of cameras have become smaller. Batteries used as power sources for the flash apparatuses also have come to be in a smaller size to have a lower voltage. As a result, it has become necessary to use a step-up circuit (a power supply device or a converter) which is made compact and arranged to stably oscillate.\nTo meet this requirement, for example, a step-up circuit disclosed in Japanese Laid-Open Patent Application No. SHO 54-102521 is arranged as shown in FIG. 6. Referring to FIG. 6, a protective resistor 102 is connected to a power supply 101. A switch 103 is connected in series to the protective resistor 102. A primary coil 104a and a feedback coil 104b of an oscillation transformer 104 are connected to the protective resistor 102 and the switch 103. The primary coil 104a is connected to the collector of an oscillation transistor 105, which is composed of a silicon transistor. The feedback coil 104b is connected to the base of the oscillation transistor 105.\nAn oscillation capacitor 106 is connected to the base circuit of the oscillation transistor 105 in parallel therewith. To a secondary coil 104c of the oscillation transformer 104 are connected a capacitor 107, a protective resistor 108 and a diode 110. A neon tube 109 is connected to the protective resistor 108.\nA main capacitor 111 is connected to the diode 110. In addition, there are provided output terminals 112a and 112b.\nWith the step-up circuit arranged in the above manner, the electric charge of the power supply 101 is supplied to the oscillation transformer 104 when the switch 103 is turned on. Then, a current flows to the base of the oscillation transistor 105, which is connected to the oscillation transformer 104. Thus, an oscillation circuit is formed by the flow of current to begin oscillation.\nThe switch 103 is arranged either to immediately turn off upon completion of its turn-on action or to turn off after the lapse of, for example, three seconds.\nThe oscillation circuit operates to continue oscillation, because energy stored in the inductance of the oscillation transformer 104 is electromagnetically induced to the feedback coil 104b and the secondary coil 104c for dissipating the energy and a current necessary for the oscillation is allowed to flow to the base circuit as the input resistance and oscillation frequency of the oscillation transistor 105 are high.\nThe signal which has been supplied to the oscillation transformer 104 and has been boosted there in the above manner is supplied to a load part which includes the main capacitor 111. The capacitor 107, which is connected to the load part, corrects the waveform of the signal by absorbing a surge voltage.\nSince the thus-corrected signal is supplied to the neon tube 109 through the protective resistor 108, the on-state of the oscillation circuit can be verified by a lighted-up state of the neon tube 109 and the off-state of the oscillation circuit can be verified by a put-out state of the neon tube 109.\nFurther, the above signal is rectified by the diode 110 and is, then, stored in the main capacitor 111. The signal stored in the main capacitor 111 is supplied from the output terminals 112a and 112b, for example, to a discharge circuit of a flash apparatus.\nThe oscillation circuit is arranged to automatically turn off in the following manner. Before the main capacitor 111, which is a load, is charged, a current of, for example, 50 mA is caused to flow to the main capacitor 111. The current comes to decrease accordingly as the main capacitor 111 is gradually charged with electric charge. When the current decreases to a value of, for example, 50 .mu.A, the base current of the base circuit also decreases accordingly. Then, the oscillation transistor 105 comes to have only a current of less than its operating point flowing there. As a result, the oscillating function of the oscillation transistor 105 comes to a stop, so that the supply of power is cut off.\nIn the conventional step-up circuit described above, the oscillation transistor 105 is connected in series to the primary winding (primary coil 104a) of the oscillation transformer 104. The base current of the oscillation transistor 105 is controlled by the charging current flowing to the main capacitor 111. The step-up circuit is thus arranged to have a current feedback system. Further, to make the start-up of the oscillation transistor 105 reliable, the feedback winding (feedback coil 104b) of the oscillation transformer 104 is also used.\nHowever, in order to charge the main capacitor 111 up to a necessary voltage level, the arrangement for using the oscillation transistor 105 for controlling the oscillation of the oscillation transformer 104 has necessitated selecting and setting, as the oscillation transistor 105, such a transistor that has a large current-amplification factor proportional to the transformer winding and has a low emitter-to-collector voltage during oscillation.\nAnother problem with the conventional step-up circuit also lies in that, depending on the setting of arrangement of the oscillation transistor 105, the size of the transistor element becomes large.\nFurther, a transistor has such an element that the current-amplification factor thereof varies with temperature. For example, at a low temperature, the current-amplification factor of the oscillation transistor becomes smaller and the oscillation comes to a stop before the charging voltage reaches a setting voltage value.\nFurther, when the load of the main capacitor becomes lighter in the last stage of a charging process, the charging current becomes smaller in proportion to the decrease of the load. Then, the base current of the oscillation transistor also becomes smaller to make the oscillation difficult. Therefore, in order to ensure a stable operation, the oscillation must be caused to be carried on by using the feedback winding. Otherwise, it is hardly possible to have a stable operation and to have the charging voltage reach the setting voltage value."}
{"text": "The invention relates to a pump cylinder a) in which a pump plunger is guided which divides the cylinder into a forward and a rearward subspace, b) which, in the forward subspace, has a connection for taking in the medium to be delivered, c) the rearward subspace of the pump cylinder having one or several electrochemical gas consumption cells for a consumption gas and d) an electric circuit being provided, which is constantly or temporarily connected with the electrodes of the gas consumption cell, as well as to a process.\nIn U.S. Patent Document U.S. Pat. No. 1,500,975, a machine is described which consists of three mutually connected cylinders and three plungers which are guided therein and are rigidly connected with one another. The three cylinders have different diameters. While the master cylinder is moved back and forth by a reciprocal admission of a vacuum and of atmospheric pressure to the spaces divided by the plunger, the two other cylinders and the plungers synchronously moved therein are used for the step-by-step compressing-on of air. In British Patent Document GB-PS 911,493, a lubricant dispenser is described. In the case of the latter, a lubricant pump presses the lubricant into a ring-shaped pipe which supplies one or several lubricating points. The ring-shaped pipe guides, by way of a pressure control valve, the excess lubricant back into the storage tank.\nSystems are known and sold, by means of which pasty masses and liquids can be continuously delivered. In one of these systems (PERMA grease dispenser of the firm G.\nSatzinger in Bad Kissingen), a cylindrical body is divided into two subspaces by a membrane acting as a plunger. One subspace contains the medium to be delivered which is supplied by way of the spout mounted at one end. The second space contains a gas generating device which continuously generates hydrogen from a corrosion element. As a result, a pressure is built up, the plunger is displaced and the medium to be delivered is pressed out of the spout.\nIn that system, the corrosion element consists of a round zinc blank with a soldered-in molybdenum rod. When activated, it is charged into a caustic potash solution. Thereby the implementation of the gas generation is defined in a constant manner. It is a function of the temperature. The delivery speed itself is also a function of the counterpressure at the spout because, when the gas generating rate is constant, the volume production of the hydrogen is inversely proportional to the pressure.\nWith respect to a short-circuit element, the cell generating the hydrogen, as described in German Patent Document DE-PS 35 32 335 A1, offers the advantage that the generating of hydrogen can be precisely controlled by way of a resistor to be controlled from the outside. This cell consists of a variant consisting of a zinc anode and of a hydrogen cathode with caustic potash solution as an electrolyte and a separator which separates the anode and the cathode. In another variant of the cell, it generates oxygen as the propellant, in which case a reducible electrode, for example, of manganese dioxide MnO2 or AgO or another metal oxide, is used as the countercathode. An oxygen precipitation electrode is used as the anode. It is characteristic of the gas generating electrodes in both cases, for oxygen as well as for hydrogen, that the electrode material has an overpotential for the gas to be generated which is as low as possible, for which in both case Raney nickel and activated carbon can be used as a catalyst.\nIn German Patent Document DE 36 21 846 A1, a delivery system is described in the case of which the liquid to be delivered is transported by way of one or a plurality of pores in a body which is lyophobic with respect to the liquid to be delivered. This measure, which divides the delivery flow into a current of discrete drops, represents a pressure lock which makes the delivery flow less dependent on the counterpressure. This pressure lock may be arranged on the oncoming flow side and on the outgoing flow side between lyophilic porous bodies and may communicate with the latter.\nSuch delivery devices can be used in multiple manners. For example, the injection of insulin solution or other medications into the body by means of such delivery devices was suggested. An earlier invention (German Patent Document DE 3711714 C2 of Apr. 7, 1987) had the object of further developing a system of the initially mentioned type such that the delivery speeds can be predetermined and controlled and that the generating of the vacuum takes place by means of devices which are as simple as possible. The essence is a driving motor which consists of an essentially rigid driving cylinder with a driving plunger movable therein or a membrane acting as a plunger. In the starting condition, this driving cylinder is filled with an electrochemically effective gas, such as hydrogen or oxygen under atmospheric pressure. Inside this driving cylinder, a gas consumption cell is situated which has a gas consumption electrode and a suitable counterelectrode, which are arranged in a cell vessel. Between the electrodes, a separator is situated which has a suitable electrolyte. The contacts of the two electrodes are accessible outside the cylinder space or are to be operated from the outside.\nWhen the closing takes place by way of an outside close-down circuit, a current flow is generated in the cell which per time unit consumes an equivalent quantity of the reaction gas and thereby generates a vacuum with respect to the outer atmosphere. This vacuum causes a movement of the driving plunger so that the latter can carry out an outside work operation.\nIn the case of the delivery system according to German Patent Document DE 3711714, the work operation consists of the movement of a pump plunger in a pump cylinder which preferably rigidly connected with the driving cylinder of the driving motor. The pump plunger is preferably also rigidly connected with the driving plunger of the driving motor. This pump cylinder contains the liquid to be delivered which, by means of the coupled movement of the two plungers is pressed out by way of the delivery spout or by way of a pressure lock.\nBased on the existing state of the art, the invention uses a different approach which has not been recognized in that it designs and utilizes the pumping cylinders known per se for collecting many different types of samples. The invention achieves this by means of the following. Accordingly, the rearward part of the pump cylinder is filled with the gas which can be electrochemically consumed, and the electric circuit is designed for causing such a current that the gas is consumed, and the pump cylinder forms a device for collecting pasty masses, liquids, gases and mobile objects for a previously definable time period.\nAlthough the present invention follows part of the teaching of German Patent Document DE 3711714 C2, in contrast to the latter, its object is not the distribution of fluid media but their collection. In many areas of manufacturing control, environmental production, medicine and biology, samples of pasty, liquid, gaseous substances or of mobile objects have to be collected, stored and tested, in which case the time periods for the collection of samples may be hours and months. This will be demonstrated by several examples:\nWaste water from many large and small enterprises is led into the public sewage system. Proof of the introduction of toxic substances can be obtained only by a continuous sampling and the analysis of the samples.\nThe collecting of gas samples in the case of conversions in technical equipment also permits the clarification and the control of reaction sequences in gases.\nThe collecting of milk samples of automatically milked cows permits the determination of the average values of an animal\"\"s fat production. Samples of the body fluids of human beings and animals are also often helpful when diagnosing diseases.\nThe collecting of small animals and microbes in forests and fields and the collection of air samples in high-pollen seasons fall into the biological field.\nIn the food production field but also when monitoring the biology of bodies of water, the proof of germs or their exclusion by collecting samples is also of interest.\nIn a simple manner, the invention now provides a device for collecting pasty masses, liquids, gases and/or other mobile objects which permits the collecting of samples for a previously definable time period.\nAn advantageous constructive variant of the invention is characterized in that the consumption gas and the consumption cell in the driving cylinder are enclosed by a deformable plastic foil.\nAccording to another preferred embodiment of the invention, the consumption gas contains oxygen and the consumption cell contains, in addition to a gas diffusion electrode for oxygen, an accumulator electrode as a counterelectrode, preferably in an alkaline electrolyte solution.\nParticularly preferably, for ensuring the secure functioning of the device, a valve is present in the take-in pipe, which valve opens in the suction direction in the case of a vacuum and has a shut-off effect in the outflow direction in the case of excess pressure.\nAccording to another particularly advantageous further development of the invention, the electric circuit is designed such that it permits a current flow in the consumption cell or cells only in the voltage window in which the gas consumption takes place and shuts off the current flow in the remaining voltage ranges. This ensures in an uncomplicated manner that the pump cylinder is operated only in the xe2x80x9ccollection and suction modexe2x80x9d.\nThe electric circuit preferably contains an optical or acoustic signal generator which generates its signal when the current flow takes place in the voltage window which permits the gas consumption.\nAccording to another advantageous variant of the invention, a starting resistor is present in addition to the adjustable working resistor, by way of which starting resistor the consumption cell will be short-circuited until the movement of the plunger starts. A manual switch or a pressure switch is also preferably provided for switching off the starting current and responds to the pressure difference at both sides of the plunger or to the movement of the plunger from the starting position of the plunger in the cylinder.\nGenerally, a consumption gas is used which is as pure as possible in order to utilize the full plunger stroke. However, by adding an inert contamination, the collecting operation can be limited because the inert gas remains in the plunger as residual gas.\nAdvantageously, time switches, photoresistors and/or photocells can be arranged as controlling components with additional direct-current cells in the electric circuit of the driving insert. As a result, the collecting operation can be limited to certain times of day, light conditions and/or other situations.\nThe invention also provides a process for collecting samples, particularly pasty masses, liquids, gases and/or other mobile objects, during a previously determinable time period in a pump cylinder characterized by filling the rearward subspace of the delivery cylinder with the consumable gas, by connecting an intake pipe in the forward cylinder space with the sample with the sample space, and by adjusting a consumption current in the gas consumption cell which, per time unit, consumes the volume of the consumption gas which is equivalent to the sample rate.\nFollowing the sampling, an analysis of the content layered corresponding to the time-related removal and/or a mixing of the content and its analysis can then take place.\nThe invention preferably provides a process for collecting mobile biological objects during a previously defined time period in a pump cylinder according to the invention, in the case of which the charging of a pilot medium into the forward cylinder space takes place before the consumption current is switched on.\nIn a further development of the invention, a medication is preferably fed to a limited body region of a living thing (for example, an organ), and, after a predetermined time period, the medication and/or other body fluids are at least partially withdrawn again from the limited body region by means of a pump cylinder according to the invention. In this case, the invention is particularly suitable for feeding to a limited body region of a living thing (for example, an organ) an overdose of a medication by means of a syringe and to withdraw this xe2x80x9coverdosexe2x80x9d again from the limited body region by means of a xe2x80x9csuction pumpxe2x80x9d according to the invention.\nThe pump cylinderxe2x80x94for example, in the form of a syringexe2x80x94is closed on top preferably by the insert which, fixed in the cylinder, may, for example be sealed in and off by an O-ring and/or be screwed to the pump cylinder. A suction insert and/or driving insert, which is also sealed off by means of an O-ring, is pushed into the insert.\nThe driving insert contains the at least one gas consumption cell or several of these cells in a series or parallel connection. In addition, there is the Zener diode in a series connection or a diode with a corresponding threshold voltage.\nThe driving insert has two positions which are characterized by a bore. This bore leads into the driving insert and permits the feeding of the consumption gas into the cylinder space, in which case the plunger is driven from its uppermost position into its lowermost position. Subsequently, the driving insert is pushed into a waiting position in which the bore of the driving insert is closed. After the activation of the driving insert, the current will flow which binds the consumption gas, pulls up the plunger and causes the collection operation to take place according to the invention.\nThe plunger may have a bore which may be closed by a frit.\nAfter the bore was utilized for charging the consumption gas, it is closed by means of a wetting liquid.\nAdditional embodiments of the invention of the invention are contained herein.\nIn the following, the invention will be explained in detail with reference to the drawing by means of embodiments."}
{"text": "In the food packaging and similar industries flexible films with barrier properties are used in large amounts to resist the transmission of water vapor and oxygen through the film. These barrier films typically have a base layer of a structurally strong material which is not very resistant to water vapor and oxygen transmission. Such films generally also include a barrier layer of a more water vapor and oxygen transmission resistant material.\nConventional barrier layers have a matrix phase of composition having an intrinsically high transmission resistance. To increase transmission resistance properties, a phase of discrete, fine particles are dispersed in the matrix phase. The dispersed phase particles are inorganic composition, such as clay, kaolin, vermiculite, silicates, and the like. It is generally understood that the dispersed phase particles are not particularly non-transmissive to water vapor and oxygen in and of themselves. However, they function to improve overall barrier film transmission resistance by providing discontinuities in the matrix and thus blocking the direct path of water and/or oxygen molecules migrating directly across the film. This causes such migrating molecules to take a tortuous, longer path which slows their flow from one side of the film to the other.\nUse of inorganic particles in barrier films has certain drawbacks. Of some concern is an adverse effect observed in the process of making barrier films with an inorganic particle dispersed phase. The matrix phase is quite often a crosslinked polymeric composition of an intrinsically transmission resistant component. Polyvinyl alcohol (PVOH) and ethylene vinyl alcohol copolymer (EVOH) are preferred polymers. Usually the inorganic particles are dispersed in a liquid coating solution that is deposited onto the base layer and dried to form the barrier layer, A preferred solvent is water. A thermally activated crosslinking agent can be included to promote the crosslink reaction of the polymers which are dissolved substantially uncrosslinked in the solvent.\nA problem arises when the inorganic particles are dispersed in the polymer matrix. A sharp rise in viscosity occurs at very low inorganic particle concentrations due to ionic charges that are inherent in inorganic mineral particles such as mica and vermiculite. This significant viscosity increase inhibits handling of the solution and the formation of a uniform and desirably thin barrier layer. High viscosity is particularly problematic when coating the wet barrier layer composition onto the base layer with gravure coating technology. There is a need to provide an effective moisture vapor and oxygen barrier transmission resistant coating and barrier film that does not suffer from excessive viscosity buildup using existing production equipment and process techniques.\nU.S. Pat. No. 7,737,200 relates to a stable, aqueous barrier composition which comprises prolamine; cold water insoluble polymer; water; water-soluble co-solvent; and stabilizer. The composition, when applied to a substrate, produces an article having a high surface energy and resistance to oil and grease penetration.\nU.S. Pat. No. 7,521,103 discloses a coating comprised of a first polymer reacted with a second polymer optionally using a crosslinker. For example, the first polymer may comprise a copolymer of polyvinyl alcohol and a vinylamine. The second polymer may comprise a copolymer of a maleic acid and an acrylic acid. The coating is present on the polymer film in an amount sufficient to increase the oxygen barrier properties of the film.\nU.S. Pat. No. 7,501,471 describes a blend comprising a poly(vinyl alcohol) stabilized vinyl acetate-ethylene polymer emulsion and a paraffin wax emulsion. A dried coating of the blend on a substrate has a hydrostatic head barrier sufficient to prevent passage of fluids but allows passage of water vapor through it. The vinyl acetate-ethylene polymer emulsion comprises a polymer containing 70 to 95 wt % vinyl acetate and 5 to 30 wt % ethylene, based on the total weight of monomers in the polymer. The paraffin wax in the paraffin wax emulsion has an average particle size of at least 0.3 microns.\nU.S. Pat. No. 7,473,729 describes an article comprising an acrylic nanocomposite film adhered to a substrate. The acrylic nanocomposite film comprises (a) a non-elastomeric acrylic polymer; and (b) an acid-treated, layered silicate filler material which has an aspect ratio of at least about 25.\nU.S. Pat. No. 7,303,797 discloses a coating layer for increasing the gas barrier performance of rigid shaped containers and films. The coating layer includes a water-borne composition comprising a water-soluble or water dispersible organic binder in the optional presence of a cross-linking agent. The organic binder includes from 10% up to less than 90% by weight, based on the dry weight of the coating layer, of an inorganic laminar mineral selected from montmorillonite, laponite, organo-modified montmorillonite and mixtures thereof.\nU.S. Pat. No. 5,981,029 describes a gas barrier film with a barrier coating. The film that provides high gas barrier ability at high humidity and under elongation, and also has durability of gas barrierability, and also with adhesiveness and post-lamination gas barrierability. The film is suggested for packaging material highly reliable in preservability. The barrier coating has inorganic stratified particles that may be provided by inorganic particles formed by overlapping one to several very thin unit crystal layers. Among them, especially, a clay is preferably used."}
{"text": "Eggs are a highly nutritious, inexpensive food commodity which is accepted by most cultures. The safety of eggs has become a global issue with emergence of the pathogen Salmonella enteritidis as a major hazard associated with the consumption of raw and semi-cooked eggs. By 1993 the occurrence of Salmonella enteritidis has reached epidemic proportions in many countries. Unlike others of the 2000 serovars of Salmonella, this organism infects the egg before the egg is laid, with the organism being transmitted to the ova or the albumen before the formation of the shell of the egg.\nIn the UK alone, there were 14500 confirmed human food poisoning cases caused by Salmonella enteritidis in 1991. In the USA, the consumption of shell eggs per annum is estimated between 50 and 65 billion eggs, of which 1 in 10000 eggs is estimated to be infected with Salmonella enteritidis resulting in a 1 in 50 chance of an individual consuming a contaminated egg in any given year. Although there has been a reduction in the occurrence of Salmonella enteritidis since 1993, due to better poultry flock management, the organism has not been eradicated from poultry flocks, especially in countries where legislative controls are not as strict.\nSalmonella enteritidis is killed by heat treatment and adequate cooking of eggs will thus kill the microorganism. However, the use of raw eggs in food as well as the consumption of partially cooked eggs is a reality, particularly in catering establishments and homes. Thus, in order to remove or reduce the risk of Salmonella enteritidis food poisoning, eggs should be pasteurized before distribution. Apart from having the effect of providing eggs substantially free of the risk of causing Salmonella enteritidis food poisoning, the shelf life of pasteurized eggs will also be significantly increased.\nThe infections of laying hens with Salmonella enteritidis, and the resultant contamination of egg contents, have resulted in a marked increase in human salmonellosis since the mid 1980's. The route of infection for Salmonella enteritidis is transovarian, whilst the route of infection for other serovars is trans-shell.\nBoth mechanical and chemical inhibitors restrict the growth of Salmonella enteritidis and other trans-shell contaminants in the albumen of fresh eggs. Fresh albumen is highly viscous and organized to confer biological structure to the egg contents, holding the vulnerable yolk away from the shell and its membranes. Fresh albumen contains factors inhibitory to the growth of microorganisms. These include lysozyme, which is an antimicrobial against Gram-positive bacteria, and conalbumin (ovotransferrin), which binds metallic ions making Fe+++ unavailable to the organisms. Additional inhibitors include ovoflavoprotein, which binds riboflavin, and anti-trypsin factors such as ovoinhibitor. As the albumen ages, it becomes less viscous and less inhibitory. The deterioration of the vitelline membrane around the yolk allows leakage of yolk contents into the albumen triggering off rapid multiplication of surviving organisms in the albumen. Temperature of egg storage is a key factor in determining the rate of albumen aging.\nThe control of Salmonella and spoilage organisms in whole shell eggs is multi-faceted. High general Salmonella infection rates in poultry and sub-clinical ovarian infection in laying hens have been attributed to such factors as intensive farming, forced moulting practices and use of contaminated animal by-products as dietary protein supplements. About 30% of eggs are sent for breakage and processing into egg products for the food processing industry. Pasteurization practices are well established for these egg products and aim to reduce the loads of Salmonella by around 5 logs. The principle source of Salmonella on eggs for breakage is fresh fecal matter on the shells and loads can be high in newly pooled fresh egg contents. Pooled yolks, pooled whole egg and pooled albumen are pasteurized at different temperatures reflecting the differences in thermal stability of the egg products.\nThe pasteurization of whole shell eggs is much more problematical than the pasteurization of egg products and must balance the required reduction of target organisms against the maintenance of albumen quality.\nTable 1 provides some information on the composition and characteristics of egg albumen and Table 2 provides more information on the characteristics of the various proteins in egg albumen:\nTABLE 1DESCRIPTION AND CHARACTERISTICS OF PROTEINS IN EGG ALBUMENAlbumenIsoelectricMolecularDenaturationProteinProteins (%)PointsWeightTemperature1Characteristicsi.Ovalbumin544.545 00084.0Phosphoglycoproteinii.Ovotransferrin126.176 00061.0Binds metallic ionsiii.Ovomucoid114.128 00079.0Inhibits trypsiniv.Ovomucin3.54.5-5.05.5-8.3 × 106—Sialoprotein, viscousv.Lysozyme (G1)3.410.714 30075.0Lyses some bacteriavi.G2 globulin4.05.53.0-4.5 × 10492.5—vii.G3A globulin4.04.8—55.0 at pH 9—viii.Ovoinhibitor1.55.149 000—Inhibits serineproteasesix.Ovoflavoprotein0.84.032 000—x.Ovomacroglobulin0.54.57.69 × 105—Strongly antigenicxi.Cystatin0.055.112 700Inhibits thiolproteasesxii.Avidin0.051063 300Binds biotin1In water or buffer\nTABLE 2DESCRIPTION OF THE CHARACTERISTICS OF EGG ALBUMEN PROTEINSProteinDescription(i) OvalbuminThis is the predominant protein in albumen. It is a phosphoglycoprotein and the onlyalbumen protein to contain free sulfhydryl group. Ovalbumin is readily denatured bymechanical stress (e.g., shaking). This is referred to as surface denaturation. Ovalbuminis thermally resistant. Ovalbumin is converted to S albumin, a more heat stable protein,during egg storage. The conversion is prevented by oiling the eggs. The denaturationtemperature of native ovalbumin and S ovalbumin are reported 84.5 and 92.5° C.(ii) ConalbuminThis is a glycoprotein that contains no phosphorus or free sulfhydryl groups. Fifteen(Ovotransferrin)disulphide bonds maintain the quaternary structure. Conalbumin is more heat sensitivethan ovalbumin but less susceptible to surface denaturation. It is reported to be stable at57° C. for 10 minutes at pH 9. Di and trivalent metallic ions are bound by conalbumin.Two atoms of Fe+++, Al+++, Cu++ and Zn++ per molecule of protein form stablecomplexes with conalbumin above pH 6. These complexes are red, colourless, yellowand colourless respectively. Complexes of conalbumin with metallic ions are resistant tothermal denaturation and proteolytic attack. The iron binding properties of conalbuminare suggested to be responsible for its anti-microbial properties.(iii) OvomucoidOvomucoid is a heat resistant glycoprotein that contains no tryptophan. Ovomucoid is atrypsin inhibitor. Chicken ovomucoid is classified as a single headed inhibitor, i.e., eachmolecule of ovomucoid binds with only one molecule of trypsin (serine proteinase).Turkey ovomucoid is double headed and duck ovomucoid triple headed. The protein isdenatured at 80° C. at pH 9.(iv) OvomucinOvomucin is a sulphated glycoprotein that may contribute to the gel like structure ofthick white in the form of flexible fibres visible under electron microscopy. The amountof ovomucin in the thick white is four times greater than in thin white. Ovomucin insolution is resistant to heat between pH 7.1 and 9.0 at 90° C. for two hours. Ovomucinand lysozyme in solution can interact to form a water-insoluble complex. Much of theinterest in ovomucin has arisen from its possible roles in maintenance of the gelstructure of thick egg white and in the process of egg white thinning. Contradictoryhypotheses have been advanced for the egg white thinning mechanisms.(v) LysozymeLysozyme is an albumin enzyme that has a lytic action on bacterial cell walls,particularly Gram positives. Lysozyme hydrolyses the Beta (1-4) linkages between N-acetylneuraminic acid and N-acetylglucosamine. In egg white, the enzyme is 50 timesmore heat sensitive than in phosphate buffer. In egg white heated to 63° C. for 10minutes, lysozyme is inactivated as the pH rises above pH 7. It has been suggested thatits role in egg white may be more important in maintaining gel structure than as an anti-microbial.(vi) OvoglobulinsThe globulin fraction consists of three proteins, G1, G2 and G3. G1 is lysozyme. G1 andG3 are glycoproteins with excellent foaming properties. G3 is also the most heatsensitive of all the albumen proteins and foaming properties of pasteurized eggs is acritical quality parameter.(vii) OvoinhibitorOvoinhibitor is a proteolytic enzyme inhibitor capable of inhibiting trypsin andchymotrypsin as well as a variety of fungal and bacterial proteases.(viii) OvoglycoproteinThis glycoprotein is an acidic glycoprotein that comprises 1% of the egg white proteins.Little is known in respect of its functionality in egg albumen.(ix) OvoflavoproteinAll the riboflavin in egg albumen is bound in the flavoprotein in a 1:1 ratio.Ovoflavoprotein is sometimes referred to as the riboflavin binding protein. The majorfunction of the riboflavin binding protein is presumably to ensure transfer of theriboflavin from the blood serum to the albumen.(x) OvomacroglobulinOvomacroglobulin is a glycoprotein with inhibitory activity against diverse proteolyticenzymes such as trypsin, papain and thermolysin.(xi) CystatinCystatin is an inhibitor of thiol proteinases such as ficin and papain. It is highly stableagainst denaturation. Since the isolation of cystatin from chicken egg white, a “superfamily” of cystatins have been isolated and characterized from mammalian tissues. It issuggested that certain cystatins may have not only the general role of protecting cellsagainst uncontrolled activity of their own proteinases, but may have a more specificaction against virus infections. This has led to research on cystatins, including eggcystatin, as possible agents for anti-viral chemotherapy.(xii) AvidinAvidin is a glycoprotein that combines with biotin to form a stable complex incapable ofabsorption by the intestinal tract of animals. In addition to its physiological importanceas a possible anti-nutrient due to its strong affinity for biotin, avidin has also beensuggested to play a role as an anti-microbial. Inhibition of microbial growth may arisenot only from the unavailability of avidin-bound biotin, but also from the ability to bindvarious Gram-negative and Gram positive bacteria.\nTable 1 is only a generalization as there are many factors that will influence the rate and temperature at which a specific protein fraction will denature. One of the most important factors is the pH of the medium and the ionic strength. Significant differences are found between published values mainly due to differences in the pH or the ionic strength of the medium in which the protein fraction was heated. Of particular importance is the fact that globulin G3A becomes much more sensitive to denaturation than ovotransferrin at a pH above 9.0. This is significant as the ovotransferrin is generally regarded as the most heat sensitive protein. Whole shell eggs of one to two days old usually have pH ranges of more than 9.0, which create a problem if ovotransferrin is used as the indicator protein for heat damage. During the pasteurization of liquid eggs, the pH and the ionic strength is usually changed by the addition of salts or buffers to create a more stable environment for the egg proteins. For example, it is known that at a pH of 7 the globulin G3A only starts denaturing at 60° C., while at a pH of 9, almost 50% of the globulin G3A is denatured when the contents have reached 60° C. Globulin G3A plays a major role in the turbidity that is first observed at the initial stages of heat treatment of eggs. The egg albumen is still liquid, but develops a “milky” or turbid appearance as the globulin G3A denatures starting at 55° C. Denaturation of the globulin G3A is immediate at 60° C. with 50% denaturation after 5 minutes exposure. As the target organism Salmonella enteritidis needs to be destroyed without damage to the egg proteins, it becomes clear that efficient heat transfer to the eggs is necessary in order to give maximum heat exposure to the microorganisms, but minimum heat exposure to the egg proteins.\nEgg yolk proteins are generally more heat stable than albumen proteins. However, as the target microorganisms are inside the yolk, yolk temperature has to be raised high enough without damage to the albumen. This is a major problem with conventional systems or processes of which the Applicant is aware using water or steam as heating media from the outside of the eggs. By the time the yolk has reached critical temperatures to kill the Salmonella enteritidis, the albumen had been exposed for extensive periods to the heat, causing the occurrence of excessive amounts of precipitate giving the albumen a milky colour. As egg albumen is generally more heat sensitive than egg yolk, egg albumen should be used as the indicator protein for egg quality assessment after processing.\nAlbumen or egg white consists of four distinct layers: outer thin white, viscous or thick white, inner thin white and the chalaziferous layer. The total solids content is 11 to 13%. Protein is the major constituent with up to 1% carbohydrate. Free carbohydrate is around 0.5% as glucose, with the rest as glycoprotein (mannose and galactose). Table 3 provides information on the layer structure and moisture content of albumen:\nTABLE 3LAYER STRUCTURE AND MOISTURE CONTENT OF ALBUMEN% of albumenLayerMeanRange% MoistureOuter thin white23.210-6088.8Thick white57.330 to 8087.6Inner thin white16.8 1-4086.4Chalaziferous2.784.3(including chalaza)\nTable 4 provides some information on the composition of albumen, yolk and whole egg:\nTABLE 4LAYER STRUCTURE AND MOISTURE CONTENT OF ALBUMENEgg ComponentProtein (%)Lipid (%)Carbohydrate (%)Ash (%)Albumen 9.7-10.60.030.4-0.90.5-0.6Yolk15.7-16.631.8-35.50.2-1.01.1Whole egg12.8-13.410.5-11.60.3-1.00.8-1.0"}
{"text": "1). Field of the Invention\nThe field of this invention is directed to whistles designed to be placed within the mouth of a human and used to produce a high pitched shrill sound by blowing air through the whistle.\n2). Description of the prior Art\nThe construction of whistles has long been known also the construction of whistles out of a confectionery product, such as hard candy, has long been known. However, heretofore, whistles constructed of the prior art have not been designed in such a way as to maximize the creation of a high pitched shrill sound.\nWhistles are commonly used as a self entertainment device. However, whistles do have utility in that they may be used as a means of getting someone's attention located at some distance away from the whistler. Also, whistles can be used as a means of instructing animals to do certain things, such as calling cows to a barn.\nA primary objective of the present invention is to construct a whistle of superior quality than previously constructed whistles that are able to produce an extremely loud high pitched shrill sound.\nAnother objective of the present invention is to construct a whistle which can be readily manufactured from either a permanent type of material, such as plastic, or can be manufactured from an edible material such as a hard candy.\nAnother objective of the present invention is to construct a whistle which can be manufactured inexpensively and thereby sold to the ultimate consumer at an inexpensive price.\nThe whistle of the present invention is constructed of a body which has a top wall and a bottom wall which are interconnected by a sidewall. An air entry through hole is formed between the top wall and the bottom wall with this air entry through hole being located directly adjacent the sidewall. Formed within the body is an air exhaust chamber which is tapered and has straight walls connected between a narrow tip and a widened base. The air entry through hole connects with the narrowed tip. The widened base forms an air exhaust opening within the top wall. The air exhaust chamber is slantingly located within the body permitting a support stick to be mounted within the body. The portion of the body including the air entry hole is to be located within a mouth of a user, and by forcing air into the air entry opening and into the air exhaust chamber a high pitched shrill sound is produced."}
{"text": "A social network service may be a computer or web-based service that enables users to establish links or connections with persons for the purpose of sharing information with one another. Some social network services aim to enable friends and family to communicate and share with one another, while others are specifically directed to business users with a goal of facilitating the establishment of professional networks and the sharing of business information. For purposes of the present disclosure, the terms “social network” and “social network service” are used in a broad sense and are meant to encompass services aimed at connecting friends and family (often referred to simply as “social networks”), as well as services that are specifically directed to enabling business people to connect and share business information (also commonly referred to as “social networks” but sometimes referred to as “business networks” or “professional networks”)."}
{"text": "The present invention relates to a hydraulic actuator preferably used for controlling braking fluid of an anti-lock braking system (hereinafter, abbreviated by ABS) installed in an automotive vehicle.\nIn the anti-lock braking system, the hydraulic actuator is operated to precisely and promptly increase and decrease the hydraulic pressure of each wheel cylinder in accordance with rotational conditions of each wheel to prevent the wheel from locking during the braking operation.\nAn object of the present invention is to secure a sufficient damper capacity for an ABS actuator to reduce the noise caused during an ABS operation.\nTo accomplish the above and other related objects, the present invention provides a first ABS actuator comprising a housing having a first predetermined surface where a bore is formed. A damper, defined in the bore, has an axis extending in a depth direction of the bore. A cup-shaped diaphragm, disposed in the bore, divides an inside space of the bore into a first chamber and a second chamber. The cup-shaped diaphragm has an opening at one axial end and a bottom at an opposed axial end. The opening of the cup-shaped diaphragm becomes a leading side when the cup-shaped diaphragm is inserted into the bore. A plug-like wall member closes an opening of the bore. A damper chamber is constituted by the first chamber and includes an inside space of the cup-shaped diaphragm. The housing has a second predetermined surface where a control valve is installed. The second predetermined surface is substantially perpendicular to the first predetermined surface. The housing has a third predetermined surface where a first port is formed. The third predetermined surface is substantially perpendicular to the first predetermined surface and is opposed to the second predetermined surface. And, a main conduit line connects the first port to the control valve.\nIn this first ABS actuator, the plug-like wall member has a protruding portion for pushing the bottom of the cup-shaped diaphragm. The protruding portion is spaced from a cylindrical wall of the bore. An annular space is formed between the protruding portion and the cylindrical wall of the bore. The annular space constitutes the second chamber. And, a first conduit line has its inlet opened to the annular space provided around the protruding portion and extends from the annular space as part of the main conduit line.\nWith this arrangement, the damper diameter can be increased to a level equivalent to the diameter of the bore. This realizes a large damper capacity. Thus, it becomes possible to reduce the noise caused during the ABS operation.\nFurthermore, the present invention provides a second ABS actuator comprising a housing having a first predetermined surface where a bore is formed. A damper, defined in the bore, has an axis extending in a depth direction of the bore. A cup-shaped diaphragm, disposed in the bore, divides an inside space of the bore into a first chamber and a second chamber. The cup-shaped diaphragm has an opening at one axial end and a bottom at an opposed axial end. The opening of the cup-shaped diaphragm becomes a leading side when the cup-shaped diaphragm is inserted into the bore. A plug-like wall member closes an opening of the bore. A damper chamber is constituted by the first chamber and includes an inside space of the cup-shaped diaphragm. The housing has a second predetermined surface where a control valve is installed. The second predetermined surface is substantially perpendicular to the first predetermined surface. The housing has a third predetermined surface where a first port is formed. The third predetermined surface is substantially perpendicular to the first predetermined surface and is opposed to the second predetermined surface. A second port is formed on the first predetermined surface where the damper is installed. And, a main conduit line is formed in the housing so as to communicate with the second chamber. The main conduit line connects the first port to the second port via the control valve.\nIn this second ABS actuator, the first port is offset from the control valve in an axial direction of the damper. The control valve is located far from the first predetermined surface compared with the first port. And, the main conduit line connects the second chamber to the control valve via a first conduit line extending obliquely with respect to the axis of the damper from a cylindrical wall of the bore.\nIn this manner, providing the first conduit line extending obliquely with respect to the axis of the damper from the cylindrical wall of the bore makes it possible to connect the second chamber to the control valve under a condition where the first port is offset from the control valve in the axial direction of the damper and the control valve is located far from the first predetermined surface compared with the first port. Accordingly, in addition to the above-described effect of the first ABS actuator, it becomes possible to downsize the housing in the axial direction of the damper.\nIn the above-described second ABS actuator, It is preferable that the plug-like wall member has a protruding portion for pushing the bottom of the cup-shaped diaphragm. And, the protruding portion is spaced from the cylindrical wall of the bore. An annular space is formed between the protruding portion and the cylindrical wall of the bore constitutes the second chamber. The first conduit line has its inlet opened to the annular space provided around the protruding portion. With this arrangement, it becomes possible to fix the position of the cup-shaped diaphragm by the protruding portion.\nAccording to a preferred embodiment of the present invention, the protruding portion of the plug-like wall member is a circular protruding portion formed at an inner end of the plug-like wall member. A diameter of the circular protruding portion is smaller than that of a main body of the plug-like wall member. The circular protruding portion presses the bottom of the cup-shaped diaphragm when the plug-like wall member is inserted into the bore, so that an insertion amount of the circular protruding portion substantially determines an axial position of the cup-shaped diaphragm in the bore.\nPreferably, the cylindrical protruding portion is spaced from the cylindrical wall of the bore when the plug-like wall member is inserted in the bore so as to close the opening of the bore, so that the second chamber surrounds the cylindrical protruding portion of the plug-like wall member.\nAccording to the preferred embodiment of the present invention, the main conduit line further comprises a second conduit line which extends from the control valve in a direction substantially perpendicular to the second predetermined surface and is connected to the first conduit line.\nThe pressurizing control valve is offset from the annular space provided around the protruding portion of the plug-like wall member in a direction normal to the first predetermined surface. The first conduit line extends straight and obliquely from the annular space and is connected to the second conduit line.\nFor example, two control valves are installed adjacent to each other on the second predetermined surface. One of the two control valves is directly connected to the first conduit line via the second conduit line. And, the other of the two control valves is connected to the first conduit line via the second conduit line and a third conduit line connecting the two control valves.\nIn this manner, when two control valves are disposed adjacent to each other, it is desirable that one of two control valves is directly connected to the first conduit line via the second conduit line. The other control valve can be connected to the first conduit line via the second and third conduit lines.\nFurthermore, it is preferable that a through-hole is opened on the bottom of the cup-shaped diaphragm. The through-hole serves as an orifice connecting the first chamber to the second chamber.\nIn this case, it is further preferable that a region where the through-hole is opened is thin in thickness compared with the remaining portion of the bottom of the cup-shaped diaphragm. For example, a recessed region is formed on the bottom so as to surround the through-hole. The recessed region has a sufficient depth so that a burr formed around the through-hole remains in a retracted position with respect to a surface of the bottom.\nAlthough the burr is inherently produced along an annular periphery of a through-hole during a formation of the through-hole, the above-described arrangement prevents the burr from protruding from the surface of the bottom. Thus, it becomes possible to prevent the burr from peeling off the cup-shaped diaphragm.\nFurthermore, it is preferable that the cup-shaped diaphragm has a press-fitting portion partly formed on a cylindrical wall thereof. The press-fitting portion has an enlarged diameter larger than the diameter of other portion of the cup-shaped diaphragm. And, the cup-shaped diaphragm is firmly fixed in the bore by the press-fitting portion.\nPreferably, the press-fitting portion protrudes radially outward from a surface of the cylindrical wall of the cup-shaped diaphragm. The press-fitting portion extends in a circumferential direction of the cylindrical wall of the cup-shaped diaphragm.\nIn this manner, forming the press-fitting portion on the cylindrical wall of the cup-shaped diaphragm makes it possible to firmly fix the cup-shaped diaphragm in the bore.\nMoreover, it is preferable that an opening periphery of the cup-shaped diaphragm is squeezed in diameter compared with other portion of the cup-shaped diaphragm.\nPreferably, the opening periphery of the cup-shaped diaphragm is curled radially inward so as to have a smaller diameter compared with a cylindrical wall of the cup-shaped diaphragm, so that the cup-shaped diaphragm can be easily inserted into the cylindrical bore from the opening periphery of the cup-shaped diaphragm.\nIn this manner, squeezing the diameter of the opening periphery of the cup-shaped diaphragm makes it possible to easily insert the cup-shaped diaphragm in the bore.\nPreferably, the cup-shaped diaphragm has a diameter substantially identical with that of the bore. The cup-shaped diaphragm has an axis identical with that of the bore."}
{"text": "In a semiconductor device manufacturing process, various chemical liquids such as, for example, a resist, is supplied to a semiconductor wafer (hereinafter, referred to as a “wafer”) as a substrate, using a coating apparatus that is one of liquid processing apparatuses. The chemical liquid, scattered or overflowing from the wafer, flows to the bottom portion of, for example, a cup that forms the coating apparatus and is installed to surround the wafer. Then, the chemical liquid flows, as waste liquid, in a liquid discharge pipe connected to the bottom portion of the cup to be removed. The liquid discharge pipe is installed to be inclined, for example, downwardly so that the chemical liquid naturally flows by weight. Japanese Patent Laid-Open Publication No. 2012-33886 discloses a coating apparatus the cup and the liquid discharge pipe as described above.\nIn some cases, a chemical liquid having a relatively high viscosity may be used as the chemical liquid. In order to ensure that the chemical liquid flows in the liquid discharge pipe even if the viscosity of the chemical liquid is so high, the liquid discharge pipe is installed at a relatively large angle with respect to a horizontal plane. However, when the liquid discharge pipe is installed in this way, the height occupied by the liquid discharge pipe at the lower side of the cup is increased, which makes it difficult to miniaturize the coating apparatus.\nWhile descriptions have been made on the waste liquid from the cup, in some cases, the coating apparatus may supply a chemical liquid to a place other than the cup, and the chemical liquid becomes the waste liquid, as will be described in detail in the description of exemplary embodiments of the present disclosure. A liquid discharge pipe configured to allow the waste liquid to flow therein also has an increased height due to a reason that is the same as the liquid discharge pipe connected to the cup, which may disturb miniaturization of the coating apparatus. The problems encountered in discharging the chemical liquid having a high viscosity as described above is not limited to an apparatus that performs a processing on a wafer like the aforementioned coating apparatus, i.e. an apparatus that performs a so-called pre-processing of a semiconductor manufacturing process. In a post-processing of the semiconductor manufacturing process, a wafer is cut to form a chip, then a chemical liquid such as, for example, a liquid resin is supplied to the chip, and the chemical liquid is cured to form a package that covers the chip. In general, the chemical liquid for forming a package also has a high viscosity. Therefore, when a liquid discharge pipe for discharging the chemical liquid is installed, the height may be increased like the liquid discharge pipe that is connected to the cup as described above.\nJapanese Patent Laid-Open Publication No. 2012-33886 discloses a system that is provided with the liquid discharge path as described above and a tank that collects waste liquid, but does not disclose a means for solving the problems as described above. Japanese Patent Laid-Open Publication No. H5-7674 discloses a technique for mixing mixing-target materials within a tank by forming vortex flows within the tank, but does not disclose the problems as described above. The mixing target materials flow in a pipe connected to the lower side of the tank and is discharged from the tank by opening/closing a valve installed in the pipe, which requires burdens in opening/closing the valve and manufacturing costs due to the installation of the valve."}
{"text": "With the advancement of transportations and the proliferation of trade barriers among nations, goods from one country can be easily transported and sold to another country. Consumers worldwide benefit the most from this advancement and proliferation as foreign goods are no longer out-of-reach due to either the mere impracticality of shipping as well as exorbitant prices. Despite these developments, international or domestic shipping cost is still an inevitable overhead expense for goods sold in any country.\nShipping cost is determined by both weight and size. To reduce weight, there must be innovative designs that would minimize materials used yet provide proper structural integrity to render a product suitable for its intended use. To reduce volume, parts must be innovatively designed to occupy the smallest footprint possible for ease of packaging and transportation. Designing a product with shipping economy and ease of distribution in mind, the product would have a much higher market penetration power than a similar product not so designed.\nThe present invention is a rectangular shape flame heater and a pyramid shape flame heater. Both products are rather tall and bulky when fully assembled. The nature of a heater demands strong rigidity and high durability per safety and reliability expectations. For these reasons, traditional flame heaters are made by welding solid metallic parts together. They are therefore very heavy and bulky for shipping, distribution and relocation purposes.\nTo overcome these shortcomings, the present invention provides a knock-down design of a rectangular shape flame heater and a knock-down design of a pyramid shape flame heater. Before being fully assembled, parts of either heater occupy a space volume many times smaller than the space volume of a fully assembled heater. The pre-assembled compact size is ideal for economy of shipment and ease of distribution."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical component such as a polarizing beam splitter or a phase plate in which a multilayer film is provided on a surface of a substrate to polarize incident light, and in particular, to an optical component having a multilayer film with a high aspect ratio as well as a method of manufacturing the same.\n2. Related Background Art\nHitherto, an optical component such as a lens, used in an electronic equipment is based on an optical interference and has thin layers stacked on a surface of an optical element as an antireflection film. Further, on the basis of a similar optical interference, thin layers are stacked on a surface of a substrate composed of a transparent optical element to polarize an incident light, thereby forming a polarizing beam splitter or a phase plate. Such an interference film is inexpensive and has good reproducibility, and is thus put to practical use widely.\nFIG. 5 shows an example of a conventionally used interference film cubic type polarizing beam splitter. In FIG. 5, reference numerals 12a and 12b denote transparent body substrates each having a surface inclined at 45°. Reference numerals 13 and 14 denote a low refractive index layer and a high refractive index layer, respectively. Further, the arrows shown in the figure indicate the direction in which a light travels. A normal light has a S polarization beam that is a vibration component having a light vibrating direction perpendicular to the drawing of FIG. 5 and a P polarization beam that is a vibration component having a light vibrating direction parallel to the drawing of FIG. 5. The polarizing beam splitter can split an incident light into two components as shown in the figure by allowing the P polarization beam to pass therethrough while reflecting the S polarization beam.\nHowever, although the cubic type polarizing beam splitter shown in FIG. 5 is very easy to manufacture and has good reproducibility, there is a problem as to its optical characteristics that the reflection and transmission characteristics change markedly depending on the incident angle of light.\nFor example, when the size of a device such as a liquid crystal projector is to be reduced, its optical path length must be shortened. The use of the conventionally used interference film cubic type polarizing beam splitter enables the optical path length to be shortened. However, because the angle at which a light is incident on the polarizing beam splitter varies markedly depending on the position of the incidence, the reflection or transmission characteristics are not constant, thus making it impossible to obtain uniform spectral characteristics within a plane. Therefore, to further improve the functions of electronic equipments, there is a need for an optical component such as a polarizing beam splitter or a phase plate that provides good optical characteristics, regardless of the incident angle of light.\nTo meet this need, a polarizer has recently been proposed which polarizes an incident light utilizing the anisotropy of propagation characteristics of the periodic structure of a transparent body and in which high refractive index layers and low refractive index layers are repeatedly stacked alternately on each other to provide a groove form that is bent periodically at a pitch less than the wavelength of visible light as disclosed in Japanese Patent Application Laid-Open No. 2000-56133. This publication describes the use of Si, GaAs, TiO2, and TaO2 as the high refractive index layer and the use of SiO2 as the low refractive index layer. Further, this publication describes an example in which the high refractive index and the low refractive index layers has a periodic thickness of 0.32 μm and in which the groove has a pitch of 0.4 μm and a depth of 0.2 μm.\nWith the configuration described in Japanese Patent Application Laid-Open No. 2000-56133 above, the polarizer can effectively function as a polarizing beam splitter for a light of a wavelength of 1 μm. However, the polarizer is difficult to use if an incident light to be polarized is a visible light having a wavelength of 400 to 700 nm. That is, Si and GaAs that constitute the high refractive index layer has a low transmissivity for visible light. Thus, the polarizer is difficult to use in an optical component such as a liquid crystal projector which must allow a light of a visible range to pass therethrough. Further, because the difference between the refractive index of TiO2 and TaO2 constituting the high refractive index layer and that of SiO2 constituting the low refractive index layer is too small, sufficient polarization characteristics cannot be obtained.\nOPTICS LETTERS (Rong-Chung et al., Vol. 21, No. 10, p. 761, 1996) describes a polarizing beam splitter comprising a plurality of dielectric multilayer films arranged in lines on a transparent substrate as shown in FIG. 6. In FIG. 6, reference numeral 21 denotes an optical component such as a beam splitter or a polarizing plate. Reference numeral 22 denotes a transparent substrate as an optical element, and reference numerals 23 and 24 denote low refractive index and high refractive index layers, respectively. The low refractive index layers 23 and the high refractive index layers 24 are alternately stacked on a surface of the transparent substrate 22 to form an HL alternate layer 25. Reference numerals 26a, 26b, 26c, 26d denote dielectric multilayer films each formed in a line. It is known that the line-shaped dielectric multilayer films 26a, 26b, 26c, 26d which form periodic recesses and protrusions function as a diffraction grating and have excellent characteristics as an optical component that polarizes an incident light. Specifically, even if an incident light is a visible light having a wavelength of 400 to 700 nm, good polarization characteristics can in principle be obtained by making the line width of each of the line-shaped dielectric multilayer films less than the wavelength of the visible light, preferably 0.1 μm or less. Further, for its manufacturing method, a patterning technology for a super LSI can be applied to manufacture the periodic recesses and protrusions using dielectric multilayer films. Examples of known patterning technology include wet etching using an etchant such as hydrofluoric acid, dry etching such as reactive ion etching, or the like.\nHowever, when a polarizing beam splitter is manufactured using the line-shaped dielectric multilayer films shown in FIG. 6, a wet process will be performed in an etching step or washing step, as described above. Further, a drying step is required after this wet process. Thus, a processing liquid flowing into the recesses during the above steps pulls the protrusions by a capillary force, so that a stress is applied to the dielectric multiplayer films as the protrusions to incline them. This tendency is significant when in the recesses and protrusions formed by the line-shaped dielectric multilayer films, the thickness of the dielectric multilayer film is relatively greater as compared with the width of the recess, i.e., the aspect ratio of the recesses and protrusions is high. This tendency becomes very significant particularly when the aspect ratio is 4 or more.\nAs a manufacturing method for avoiding this, Japanese Patent Application Laid-Open No. 2001-165568 discloses a drying method utilizing a supercritical state. According to this method, a processed pattern is immersed into normal hexane as a nonpolar solvent, and then sealed together with liquid carbon dioxide in a reaction chamber, so that the normal hexane is replaced with the liquid carbon dioxide to dry the processed pattern. This method is effective in preventing defects in the pattern. However, the inner pressure of the reaction chamber must be set to 7.5 MPa, thus requiring a high pressure container. Further, an organic solvent is required as a nonpolar solvent, which may pollute the environment."}
{"text": "1. Field of the Invention\nThe present invention relates to a cathode ray tube, and more particularly, to a neck of a funnel and a stem sealed to the neck of a cathode ray tube.\n2. Description of the Related Art\nA cathode ray tube includes a phosphor layer on the inside thereof, a panel to which a shadow mask spaced from the phosphor layer is secured, a funnel having a neck and a cone portion connected to the panel, an electron gun housed in the neck for emitting an electron beam, a deflection yoke fixedly installed around the cone portion, and a stem sealed to one end of the neck for mounting the electron gun.\nThe cathode ray tube operates as follows. First, if a heater installed on the inside of a cathode of the electron gun generates heat, electrons are emitted from oxide coated on the top of the cathode. Then, the electrons emitted from the cathode pass through each electrode arranged at regular intervals to form the electron beam of the desired characteristics. The formed electron beam is deflected by a magnetic field produced by the deflection yoke, passes through the shadow mask, and collides with the phosphor layer coated on the inside of the panel to light up phosphors, thereby creating a screenful of an image.\nTo smoothly perform the operation as described above, the interior of the cathode ray tube must maintain a vacuum. To accomplish this, a stem having a plurality of stem pins for supplying voltages to the electron gun and an exhaust pipe for exhaustion is introduced into one end of the neck, and a portion at which the side of a stem flange contacts the inside of the neck is fused and sealed off. The interior of the cathode ray tube is evacuated through the exhaust pipe to a vacuum and the exhaust pipe is then fused and sealed.\nThe stem includes the stem flange formed in the shape of a flat disk whose diameter is smaller than the inside diameter of a sealing portion, the plurality of stem pins arranged in a round shape to pass through the stem flange for introducing signal voltages from an external circuit, and a plurality of stem mounds convexly built of glass, which is the material of the stem, for holding the plurality of stem pins and preventing the loss of vacuum, and the exhaust pipe formed in the central part of the stem flange for evacuating cathode ray tube to a vacuum. Here, the diameter of an inner stem pin circle of interior stem pins connected to the electrodes of the electron gun is equal to that of exterior stem pins connected to a socket for applying a predetermined voltage of each electrode of the electron gun.\nAs described above, the exterior stem pins are combined with the sockets installed in a chassis. For example, the stem used in the neck having a diameter of 22.5 mm (millimeters) is fit into a socket for 22.5 mm, and the stem used in the neck having a diameter of 29.1 mm is fit into a socket for 29.1 mm. However, this raises a problem in that a cathode ray tube having a neck of diameter 29.1 mm are not compatible with that having a neck of a diameter 22.5 mm since chassis for 29.1 mm has been chiefly manufactured in a market for monitors of 15 or more inches (diagonal measurement of screen).\nRecently, an electric potential applied to a focusing electrode of an electron gun tends to increase due to a flat panel of a cathode ray tube and increased dynamic focusing modulation. Furthermore, current must be applied to coils of a deflection yoke to deflect electron beams emitted from the electron gun in a cathode ray tube. Since a smaller amount of current is consumed as the diameter of a neck decreases, the diameter of the neck tends to be less for low power consumption.\nHowever, high electric potential and small diameter of a neck results in large spherical aberrations due to a decreased size of electrodes of an electron gun. To prevent this, the number of electrodes of an electron gun should be increased. Since the increased number of electrodes increases the number of stem pins accordingly, problems associated with a breakdown voltage may occur. To solve the breakdown voltage problems, one empty pin may be inserted on either side of a high voltage stem pin. However, since insertion of empty pins may result in the increased number of stem pins, a stem having a large stem pin circle is required.\nFurthermore, to solve the breakdown voltage problems, the diameter of a neck may be made larger, and the diameter of a stem flange may be made larger to seal it to one end of the neck. However, the large diameter of the neck results in high power consumption and sealing the stem flange to the neck end may require an extra device and drop a yield rate.\nIt is more difficult to fuse and seal the stem having a large stem pin circle to a low deflection cathode ray tube having a narrow neck of a diameter 22.5 mm than to a cathode ray tube having a neck of a diameter 29.1 mm presently widely used. Furthermore, if a stem mount is formed very close to a connecting portion where the stem flange is fused and sealed to the end of the neck, cracks may occur easily at the connecting portion of the stem flange and the neck.\nA neck of a cathode ray tube is disclosed in U.S. Pat. No. 6,078,134 issued to Nose et al. for Narrow-neck CRT having a Large Stem Pin Circle.\nIt is therefore an object of the present invention to provide a cathode ray tube having a stem compatible with a cathode ray tube having a neck of a different diameter, which simplifies a fabrication process and increases a yield rate by sealing the compatible stem to a neck with an existing sealing device.\nIt is another object of the present invention to provide a cathode ray tube having a stem-sealing region of a neck, where the inside diameter of the sealing region is increased so that it is easier to fuse and seal a stem having a large stem pin circle to the narrow neck of the cathode ray tube with low deflection.\nAccordingly, to achieve the above and other objects, the present invention provides a cathode ray tube. The cathode ray tube includes a panel in which a phosphor layer is formed, a funnel connected to the panel, the funnel including a neck having a region for housing an electron gun and a region to which a stem is sealed, and a stem having a plurality of stem pins, each stem pin being supported by each stem mound for applying voltage to each electrode of the electron gun. The inside diameter of the stem sealing region of the neck is greater than that of the electron gun-housing region, the diameter of an inner stem pin circle formed by interior stem pins disposed on the inside of the neck is less than that of an outer stem pin circle formed by exterior stem pins disposed on the outside thereof, a horizontal length between an outer edge of the stem mound and an interior of the neck is in the range greater than or equal to 1.0 mm and less than or equal to 2.0 mm.\nIn another embodiment, a cathode ray tube includes a panel in which a phosphor layer is formed, a funnel connected to and tapered from the panel, and a neck connected to the funnel and including an electron gun-housing region and a stem sealing region, to which a stem having a plurality of stem pins arranged in an annular shape and passing therethrough for introducing signal voltages from an external circuit is sealed, wherein D1 is 22.5xc2x10.7 mm and D2 is in the range greater than D1 and less than or equal to 24.0 mm where the outside diameters of the electron gun-housing region and the stem sealing region are D1 and D2, respectively.\nAs described above, the cathode ray tube according to an embodiment of the present invention forms an outer stem pin circle greater than an inner stem pin circle, thereby achieving compatibility with cathode ray tubes having a neck of a different diameter. The stem is sealed to the inside of the neck by making the inside diameter of the sealing region of the neck larger than that of an electron gun-housing region, thereby increasing a yield rate without the need for a special device which is otherwise required for sealing a stem to one end of a neck. Furthermore, the stem is sealed at the stem-sealing region of the neck, thereby removing glass residues or foreign material and increasing a breakdown voltage.\nThe cathode ray tube according to another embodiment of this invention increases the outside diameter of a stem-sealing region of a narrow neck so that it is easier to fuse and seal the stem having a large stem pin circle to the narrow neck of the low deflection cathode ray tube. Furthermore, the cathode ray tube increases a distance between a stem mound on a stem flange and a connecting portion of the neck, thereby preventing occurrences of crack at the connecting portion."}
{"text": "The present invention relates to dynamic random access memory (DRAM), and more specifically, to a method and apparatus for controlling data transfers to and from a dynamic random access memory.\nDynamic random access memory (DRAM) components, such as those illustrated in FIG. 1A, provide an inexpensive solid-state storage technology for today\"\"s computer systems. Digital information is maintained in the form of a charge stored on a two-dimensional array of capacitors. One such capacitor is illustrated in FIG. 1B.\nFIG. 2 illustrates a prior art memory system including DRAM with the corresponding control, address and data wires which connect the DRAM to the processor or memory controller component. In synchronous DRAMs, a write access is initiated by transmitting a row address on the address wires and by transmitting row address strobe (RAS) signal. This causes the desired row to be sensed and loaded by the column amplifiers. The column address is transmitted on the address wires and the column address strobe (CAS) signal is transmitted along with the first word of the write data WData(a,1). The data word is then received by the DRAM and written into the column amplifiers at the specified column address. This step can be repeated xe2x80x9cnxe2x80x9d times in the currently loaded row before a new row is sensed and loaded. Before a new row is sensed, the old row must be restored back to the memory core and the bit lines of the DRAM precharged.\nFIG. 3A illustrates synchronous write timing. In the figure, a, b . . . represent a row address; 1, 2 . . . n represent a column address, WData[row, col] represents the DRAM address of data words, the row address strobe (RAS) is a control signal for initiating a sense operation, and WRlTE(CAS) initiates the write operation on the column amplifiers. In the present example, the row column address delay timing parameter is equal to two clock cycles. After the row address is asserted at the first clock cycle, column addresses and write data are asserted after the delay to write the data into the DRAM array.\nFIG. 3B illustrates synchronous read timing. A processor initiates a read access by transmitting a row address on the address wires and by transmitting the row address strobe (RAS) signal. This causes the desired row to be sensed by the column amplifiers. The column address is then transmitted on the address wire and the column address strobe (CAS) signal is transmitted. The first word of the read data RData (a,1) is then transmitted by the DRAM and received by the processor. This step can be repeated xe2x80x9cnxe2x80x9d times in the currently loaded row before a new row is sensed and loaded. Before a new row is sensed, the old row must be restored back to the memory array.\nVarious attempts have been made to improve the performance of conventional DRAMs. Such attempts have resulted in DRAM architectures that deviate in varying degrees from conventional DRAM architectures. Various alternative DRAM architectures are described in detail in NEW DRAM TECHNOLOGIES, by Steven A. Przybylski, published by MicroDesign Resources, Sebastopol, Calif. (1994). Some of those architectures are generally described below.\nThe prior art includes Extended Data-Out (EDO) memory systems. In EDO DRAMs, the output buffer is controlled by signals applied to output enable (OE) and column address stobe (CAS) control lines. In general, data remains valid at the output of an EDO DRAM longer than it does for conventional DRAMs. Because the data remains valid longer, the transfer of the data to the latch in the memory controller can be overlapped with the next column precharge. As a result, burst transfers can be performed in fewer clock cycles.\nThe prior art also includes Synchronous DRAM (SDRAM) memory systems. The interface of an SDRAM includes a multiplexed address bus and a high-speed clock The high speed clock is used to synchronize the flow of addresses, data, and control on and off the DRAM, and to facilitate pipelining of operations. All address, data and control inputs are latched on the rising edge of the clock. Outputs change after the rising edge of the clock. SDRAMs typically contain a mode register. The mode register may be loaded with values which control certain operational parameters. For example, the mode register may contain a burst length value, a burst type value, and a latency mode value. The burst length value determines the length of the data bursts that the DRAM will perform. The burst type value determines the ordering of the data sent in the bursts. Typical burst orders include sequential and sub-block ordered. The latency mode value determines the number of clock cycles between a column address and the data appearing on the data bus. The appropriate value for this time interval depends largely on the operating frequency of the SDRAM. Since the SDRAM cannot detect the operating frequency, the latency mode value is programmable by a user.\nThe prior art also includes memory systems in which data transfer operations are performed by DRAMs in response to transfer requests issued to the DRAMs by a controller. Referring to FIG. 4, it illustrates a memory system in which data transfers are made in response to transfer requests. The request packet format is designed for use on a high speed multiplexed bus for communicating between master devices, such as processors, and slave devices, such as memories. The bus carries substantially all address, data, and control information needed by the master devices for communication with the slave devices coupled to the bus. The bus architecture includes the following signal transmission lines: BusCtl, BusData [8:0], BusEnable, as well as clock signal lines and power and ground lines. These lines are connected in parallel to each device.\nThe processors communicate with the DRAMs to read and write data to the memory. The processors form request packets which are communicated to the DRAMs by transmitting the bits on predetermined transmission lines at a predetermined time sequence (i.e. at predetermined clock cycles). The bus interface of the DRAM receiver processes the information received to determine the type of memory request and the number of bytes of the operation. The DRAMs then perform the memory operation indicated by the request packet.\nFIG. 5 illustrates command control information 500 that is sent in a data transfer request according to a prior art protocol. In the illustrated example, the command control information 500 is sent over a BusCtl line and a nine-bit data bus (BusData[8:0]) in six clock cycles. The command control information 500 includes groups of bits 501, 502, 504, 506 and 508 that constitute an address, an operation code consisting of six bits 510, 512, 514, 516, 518 and 520, and groups of bits 522, 524 and 528 that specify a count. The address identified in the command control information 500 specifies the target DRAM and the beginning location within the DRAM of the data on which the operation is to be performed. The count identified in the command control information 500 specifies the amount of information on which the operation is to be performed.\nOne object of the present invention is to provide a mechanism to decouple control timing from data timing.\nAnother object of the present invention is to provide mechanisms that use minimal bandwidth to determine data timing while minimizing the latency from signaling that the data transfer should terminate to the transmission of the final data packet.\nAnother object of the present invention is to provide mechanisms for arbitrarily long data transfers following a command. This may include simultaneous provision of a new column address for each data packet transferred.\nAnother object of the present invention is to provide a means to signal simultaneously with termination of the data transfer that a precharge operation should be performed.\nAnother object of the present invention is to provide mechanisms and methods for interleaving control and data information in such a fashion that pin utilization is maximized without placing latency requirements upon the DRAM core that are difficult or expensive to satisfy.\nAnother object of the present invention is to provide a mechanism for interleaving control and data information that minimizes bandwidth consumed for signaling the beginning and ending of data transfers.\nAnother object of the present invention is to provide for devices that do not always interpret the information presented at their pins. Each command provides sufficient information that all further control information related to the command can be easily determined even in the presence of control information related to previous command transfers.\nAnother object of the present invention is to provide a mechanism for optionally sequencing a series of core operations prior to data transmission and, optionally, a final core operation after data transmission is terminated\nAnother object of the present invention is to provide a DRAM core which allows a single high current RAS operation at any one time in order to minimize the cost and complexity of the DRAM.\nAnother object of the present invention is to provide an encoding of the command such that decoding space and time is minimized and functionality is maximized.\nThe present invention provides a method and apparatus for performing data transfers within a computer system. The method includes causing a controller to transmit control information on a bus. The control information specifies a data transfer operation and a beginning location of data to be transferred. The controller determines, after transmitting the control information on the bus, a desired amount of data to be transferred in the data transfer operation. The controller transmits over the bus a terminate indication at a time that is based on the desired amount of data and a beginning time of the data transfer operation. A memory device reads the control information on the bus. The memory device performs the specified data transfer operation on data stored at the beginning location. The memory device continues to perform the specified data transfer operation until detecting the terminate indication on the bus. The memory device ceases to perform the data transfer operation at a time that is based on the time at which the terminate indication is detected.\nOther objects, features, and advantages of the present invention will be apparent from the accompanying drawings and from the detailed description that follows below."}
{"text": "Portable lights which can be manually moved and suspended about a work site to aid a user to obtain desirable lighting conditions are well known. It has been the practice to use incandescent light bulbs, suitably encased in light guards, for this purpose. Such lights are often referred to as work lights, trouble lamps, extension lights, inspection lights, and the like, and are commonly employed by mechanics and other workers who require a concentration of light in a frequently changing location. Such a work light is shown in U.S. Pat. No. 4,774,647 to Kovacik et al.\nThere are several drawbacks associated with the use of incandescent light bulbs in work lights. For example, incandescent lights use a relatively large amount of electrical energy as compared to other types of lights such as fluorescent lights and LED lights; become hot during operation; and are known to fail when exposed to rough service conditions.\nFurther, due to the relatively large amount of electrical energy consumed by incandescent lights, legislative efforts are underway to phase out the use of incandescent lights. Accordingly, incandescent work lights may become obsolete in the event incandescent light bulbs become unavailable.\nIt would be desirable to produce an LED conversion module for use in an incandescent work light, wherein the LED conversion module minimizes a consumption of electrical energy and facilitates conveying the work light from site to site in a portable fashion."}
{"text": "A. Field of the Invention\nThe present invention relates to a wet friction plate within a wet power transmitting and interrupting mechanism or clutch mechanism, where the wet friction plate is provided with a cushioning member.\nB. Description of the Background Art\nIn a wet power transmitting and interrupting mechanism such as a wet clutch and a wet brake, a wet friction facing used in a wet friction plate is made of a paper or organic material or the like.\nFor instance, a wet multiple plate clutch arranged in a transmission is used for changing direction and speed of torque transmission. The multiple plate clutch includes a plurality of drive plates and a plurality of driven plates which are arranged alternatively to each other in the axial direction. Each driven plate is provided at its opposite surfaces with wet friction facings, e.g., of an organic material fixed thereto. One problem with such devices is that the surfaces engaged by the drive and driven plates are expensive to manufacture in that they must be very finely machined to be flat or planar.\nA lockup clutch of a torque converter is a clutch device for transmitting a torque by mechanically coupling a front cover and a turbine together, and is provided for improving a fuel consumption of a vehicle. The lockup clutch is formed of, e.g., a piston which can be coupled to the front cover and a damper mechanism for coupling the piston and a member at the turbine side together. A wet friction facing made of organic material is fixed to a side of the piston opposed to the front cover.\nThe damper mechanism of the lockup clutch absorbs a torque vibration caused by vibration in combustion in an engine. However, it is impossible to operate the lockup clutch in a low speed range of the vehicle, because a torque vibration at a level, which cannot be sufficiently absorbed by the damper mechanism, is generated in this low speed range. In recent years, slip control has been utilized for further improving a fuel consumption by operating the lockup clutch in a lower speed range. The slip control is performed by pressing the piston against the front cover with a small coupling force, and thereby steadily allowing controlled slippage between the piston and the front cover. When the controlled slippage is utilized, torque is transmitted through divided paths, i.e., a mechanical transmission path (the slipping surfaces) and a hydraulic transmission path. When the controlled slippage is large, torque is mechanically transmitted at a small rate, and is hydraulically transmitted at a large rate. When the controlled slippage is small, a power is mechanically transmitted at a large rate, and is hydraulically transmitted at a small rate. The controlled slippage speed is controlled by a hydraulic pressure control device which controls a difference between hydraulic pressures at opposite sides of the piston in the torque converter.\nIn a lockup clutch and a wet multiple plate clutch in the prior art, it is difficult to keep uniform contact in circumferential and radial directions between the wet friction facing and the friction surface. In other words, partial contact is liable to occur. This often results in the following and other problems.\n(1) A shudder is liable to occur when the clutch is engaged. PA1 (2) A wet friction facing is liable to wear to a higher extent. PA1 (3) In the lockup clutch, leak of a lockup hydraulic pressure is liable to occur.\nFor reducing a partial contact, the opposed friction surface must have a high flatness, which requires a high processing accuracy and therefore increases a cost.\nThe slip control of the lockup clutch in the torque converter suffers from the following problem. In the prior art, the slip control is performed, for example, in a vehicle speed range from 30 to 48 km/h with a fourth gear position. Although it is preferable to perform the slip control in a lower vehicle speed range for improving a fuel consumption, this results in the following difficulties.\nIn a low speed range, it is necessary to change a ratio of the power transmitted through a hydraulic path to a large extent, so that it is necessary to increase the slip rotation speed to a certain extent. In order to increase the slip rotation speed to a certain extent, a difference between hydraulic pressures at opposite sides of the piston member must be extremely reduced. When the pressure difference is finely controlled in this state, the piston may suddenly move toward the front cover due to change in hydraulic pressure difference. When this occurs, the torque variation is mechanically transmitted to the output member. For the above reason, it is difficult to perform the slip control at a low speed range, e.g., under 30 km/h."}
{"text": "The instant invention relates generally to C-clamps and more specifically it relates to a double C-clamp device.\nNumerous C-clamps have been provided in the prior art that are adapted to hold various types of work articles together. For example, U.S. Pat. Nos. 4,691,907 to Yang; 4,828,241 to Yang and 4,962,918 to Yang all are illustrative of such prior art. While these units may be suitable for the particular purpose to which they address, they would not be as suitable for the purpose of the present invention as hereafter described."}
{"text": "This invention is in the field of circuitry for driving light-emitting diodes (LEDs). Embodiments are more specifically directed to LED driver circuitry in transmitters of photoplethysmography (PPG) systems.\nPhotoplethysmography (PPG) is a modern technology that has proven useful for the measurement of cardiovascular function in humans. According to this technology, fixed wavelength light from a light-emitting diode (LED) is emitted into the skin of a human subject, and is sensed by a photodiode (PD) after transmission through the skin and underlying tissue. The characteristics of the sensed light allows measurement of medical parameters such as oxygenation, pulse rate, respiratory function, and the like.\nConventional PPG sensors include the well-known pulse oximeter, for example of the type that clip-on onto the finger of the patient. Pulse oximeters typically measure the oxygen saturation of circulating blood from a comparison of the absorption of light in the dermis and subcutaneous tissue at two different wavelengths. So-called “wearable” devices such as heart rate sensors also utilize this technology, but need only measure light absorption at a single wavelength.\nFIG. 1 generically illustrates the architecture of a conventional PPG system. Transmitter 3 in this system includes LED 2, which has its anode biased by the Vdd power supply voltage and its cathode coupled to ground via LED driver 4. As mentioned above, PPG pulse oximeters will included multiple LEDs 2 (e.g., red and infrared), operated in time-multiplexed fashion. When forward-biased by LED driver 4, LED 2 emits light into the patient, for example the index finger of the patient. Receiver 7 includes photodiode 6, which has its cathode biased at the Vdd power supply voltage and its anode connected to the input of amplifier 8, and which is normally reverse-biased so that photons impinging photodiode 6 will produce a current detectable by amplifier 8. In this manner, photodiode 6 senses the extent to which the light emitted by LED 2 is transmitted through the subject. The output of amplifier 8 is forwarded to the desired processing and analysis circuitry of the PPG system to determine the desired medical measurement, such as the oxygenation of the patient's blood.\nThe ability of any PPG system to accurately and precisely measure the parameter of interest is based on the signal-to-noise ratio (SNR) of the overall system, considering both its transmitter and receiver. For example, it has been observed that an SNR of at least about 30 dB, for the PPG system as a whole including the transmission channel through patient tissue, is necessary in order to measure pulse rate to an accuracy of 1 beat per minute (bpm). A complicating factor in practice is that the system SNR depends on the perfusion index of the patient, as illustrated in FIG. 2. Perfusion index is the ratio of the AC signal due to pulsatile blood flow to the DC background level of the light signal passing through the patient's peripheral tissue, and depends largely on the health and physical condition of the patient. As such, for the PPG system of FIG. 1, the perfusion index is reflected in the amplitude of the AC pulses of the received light (e.g., as output by amplifier 8) relative to its DC level. As shown in FIG. 2, the system SNR increases at higher perfusion index values.\nConversely, if the SNR of the PPG system can be increased, the system can measure the pulse rate and blood oxygenation in a wider range of patients, particularly those of poorer health and thus lower perfusion indices. FIG. 2 shows two SNR vs. perfusion index plots 9a, 9b. Plot 9a illustrates the relationship of system (transmit—channel—receive) SNR to perfusion index for the case in which the transmitter SNR is 95 dB; as known in the art, to attain the required system SNR, the SNRs for the transmitter and receiver must both be higher than that required system SNR. At that transmit SNR, the PPG system is able to measure pulse rate to an accuracy of one bpm only for patients exhibiting a perfusion index of at least about 0.06. In contrast, plot 9b illustrates that if the transmitter is able to operate at an SNR of 110 dB, pulse rate measurements at an error of 1 bpm can be made for patients with perfusion index values as low as about 0.01. Accordingly, noise in the transmitter of the PPG system is a critical factor in covering a wide range of patients.\nSo-called “wearable” electronic devices, such as fitness monitoring devices, have recently become popular. In addition to fitness monitoring devices, wearable medical monitoring devices are being contemplated for use in healthcare, for example to monitor the recovery or progress of a patient suffering from a medical condition. As such, the use of PPG to obtain oxygenation, pulse rate, and other measurements by way of a wearable device, particularly such a device that can be worn all day, is desirable. In this context, battery life becomes of critical performance.\nIn this regard, an important electrical parameter of a transmitter in a battery-powered system, such as a wearable device, is the “headroom” of the LED driver. As well-known by those in the art, it is desirable that battery-powered systems operate at low power supply voltages to reduce power consumption and to reduce the cost of the battery itself. In transmitter 3 of FIG. 1, the voltage drop across LED 2 in its operating state is defined by its material. The headroom, shown as Vhead in FIG. 1, is the voltage required by LED driver 4 beyond the LED voltage drop. Conversely, the minimum Vdd power supply voltage is the sum of the voltage drop across LED 2 and the headroom Vhead of LED driver 4. For battery powered systems, therefore, it is desirable that the headroom Vhead required by LED driver 4 be minimized, especially considering that the output voltage from conventional batteries tends to sag over time.\nAs mentioned above, it is desirable to minimize power consumption in PPG systems, particularly those in battery-powered wearable devices intended for “all-day” use. Duty cycling of LED driver 4 in transmitter 3 is a common approach to reducing system power consumption. It is therefore desirable for LED driver 4 to exhibit fast switching, and rapid settling times, so that the “on” pulse width can be reduced as much as possible and thus minimizing power consumption.\nFIGS. 3a through 3c illustrate examples of conventional LED drivers for PPG systems. The circuit of FIG. 3a is an example of a typical LED driver circuit, such as used in a PPG system as described above. In this circuit, LED 10 has its anode at the Vdd power supply voltage and its cathode connected to the drain of rise time control n-channel MOS transistor 12. Transistor 12 has its source connected to the drain of n-channel driver transistor 14, which is connected to ground via variable resistor 16. Resistor 16 operates to control the current drawn through LED 10 and transistor 12, at a resistance typically set by a digital-to-analog converter (DAC). The gate of transistor 12 receives a control voltage from rise time controller 13, which is an adjustable circuit block that controls the conduction of transistor 12 to attain the desired rise and fall times in the turn-on and turn-off characteristics of LED 10. Amplifier 18 receives a reference voltage VREF at its non-inverting (positive) input and a feedback voltage from the source of transistor 14 at its inverting (negative) input. Output voltage VGATE from amplifier 18 is applied to the gate of driver transistor 14. According to this arrangement, amplifier 18 operates to drive the gate voltage VGATE at driver transistor 14 so that reference voltage VREF at the source node of transistor 14. Reference voltage VREF is modulated to selectively forward bias LED 10.\nThe LED driver circuit of FIG. 3a provides certain advantages in a PPG system. Specifically, variable resistor 16 tends to reduce the transmitter noise in this circuit, and the ripple exhibited by this circuit is also quite low. However, it has been observed that this arrangement is vulnerable to significant input flicker noise from amplifier 16, degrading transmitter performance. The settling time of this LED driver is also quite slow, due to the bandwidth of amplifier 18. In addition, the LED driver of FIG. 3a is not conducive to implementation in low voltage, battery-powered, applications because of its large headroom voltage, specifically the sum of the drain-to-source overdrive voltages of transistors 12 and 14 plus the voltage drop across resistor 16.\nFIG. 3b illustrates a conventional LED driver with very low headroom requirements as useful for a PPG system. In this circuit, power supply 20 applies the Vdd bias to LED 10 through inductor 22. N-channel driver transistor 26 has its source-drain path connected in parallel with LED 10 between inductor 22 and ground. The gate of transistor 26 receives the output of pulse-width modulator (PWM) 24. During the “off” pulses, transistor 26 shunts the inductor current through inductor 22 to ground; during the “on” pulses, the Vdd power supply voltage forward biases LED 10, such that the inductor current is conducted through LED 10 to ground. FIG. 3b illustrates the behavior of output current Iout through LED 10 over a sequence of pulses from PWM 24, illustrating that the output current Iout appears as a sequence of triangle waves. While the headroom required by this LED driver is quite low, it has been observed that transmitter noise is quite high in this arrangement, which reduces the patient coverage as discussed above relative to FIG. 2. In addition, significant ripple is present in the LED driver of FIG. 3b. \nIn the circuit of FIG. 3c, LED 10 has its anode at the Vdd power supply voltage and its cathode connected to the drain of n-channel MOS transistor 14; the source of transistor 14 is at ground. A reference voltage VREF is applied to the gate of transistor 14, and is modulated to turn LED 10 on and off, thus controlling the emission of light. This simple driver of FIG. 3c has a low headroom voltage of only the drain-to-source voltage overdrive of transistor 14, and exhibits no ripple. However, because the driver of FIG. 3c is quite noisy, its use in the transmitter of a PPG system will have limited patient coverage, as discussed above relative to FIG. 2.\nBy way of further background, auto-zeroing techniques for removing offset voltage and drift of operational amplifiers are known in the art, as described in Kugelstadt, “Auto-zero amplifiers ease the design of high-precision circuits”, Analog Applications Journal, 2Q 2005 (Texas Instruments Incorporated), pp. 19-28, incorporated herein by reference."}
{"text": "1. Field of the Invention\nThis invention related to bird decoys on a land surface and more specifically motorized bird decoys simulating a feeding position.\n2. Description of the Prior Art\nBird decoys are used by hunters to attract game birds. These decoys come in several configurations. Some float on the water naturally, others float on the water with an electric motor with a remote control, and others have feet to stand on solid terrain. Bird decoys also come in various sizes, from the small ducks to the large geese and turkeys.\nThere are many patents on bird decoys and some decoys are quite sophisticated and expensive with complete remote control device for water performance. Only the decoys adopted for land terrain are relevant to the present invention.\nU.S. Pat. No. 5,036,614 to Jackson simulates the feeding action of a bird. Jackson must provide a pivoting mechanism, while the present invention utilizes the design of the bird decoy for pivoting action.\nU.S. Pat. No. 3,916,553 to Lynch et al utilizes a mechanism to move the head and neck of the bird decoy for simulated feeding.\nU.S. Pat. No. 5,289,684 to Denny et al is a radio controlled bird decoy that moves the head of the bird up and down and from side to side.\nU.S. Pat. No. 5,233,780 to Overholt is a bird decoy that gives radio control auditory and visual signals, including movement.\nU.S. Pat. No. 5,459,958 to Rieinke is a bird decoy which is string activated to provide mating movements.\nU.S. Pat. No. 2,849,823 to Miller is a motor operated bird decoy that simulates shaking water off of stretching wings.\nU.S. Pat. No. 4,896,448 to Jackson described a bird decoy having a mechanism to flap wings.\nNone of the above patents enjoy utilizing an existing bird decoy and adding movement means without modification to the decoy. What is needed is a simple, lowcost mechanism that can be added to an existing bird decoy to simulate natural movement of birds. Accordingly, a fuller understanding of the invention may be obtained by referring to the Summary of the Invention, and the Detailed Description of the Preferred Embodiment, in addition to the scope of the invention defined by he Claims taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates generally to drive belts for motorized apparatus, and more particularly, to a temporary drive belt that can be installed without the use of a wrench.\n2. Description of the Prior Art\nDrive belts, such as fan belts to be installed on an automobile engine in emergency situations when a conventional fan belt breaks or is in an unsafely weakened condition, are well known in the art. For example, the following U.S. Pat. Nos. disclose emergency drive belts: 3,747,165 of Brown; 4,207,776 of Helt et al; 4,376,631 of Garza; 4,795,410 of Alderfer; and 4,861,322 of Reddick. The belt of Brown requires an extraneous razor blade to cut off excess belt material when the emergency belt is installed. The belt of Helt et al requires a plurality of connectable links to obtain a belt of the desired length. The belt of Garza employs a plurality of teeth projecting from a rigid external sleeve to dig into the two ends of the belt and retain the two ends of the belt to form an endless loop. No external force is used to assure that the teeth continue to dig into belt material under the stress of use on the engine. Also, the digging into the belt ends can cause damage to the belt ends. The belt of Alderfer requires a tubular belt and uses a barbed shank, and end of which is inserted into each of the respective ends of the tubular belt. The barbs dig into the two ends of the tubular belt. The tubular belt is hollow, and as such, contains a reduced amount of belt material, whereby the strength of the belt is diminished. Moreover, the hollow belt undergoes an excessive amount of flexing, compared with a solid belt, and the excessive flexing can contribute to excessive heat build up and premature wear. The belt of Reddick uses a simple coupling strip of material for connecting two ends of the belt. No outside force is exerted on the coupling strip. As a result, the strength of the connection between the two ends of the belt may be overcome by the stresses of an operating engine, especially when the engine is started from a non-running condition.\nThus, while the foregoing body of prior art indicates it to be well known to use emergency drive belts generally, the provision of a simple, cost effective, easily installed, and secure device is not contemplated. The prior art described above teaches the use of an extraneous razor blade to cut off excess belt material when the emergency belt is installed. Also, the prior art discloses a plurality of connectable links to obtain a belt of the desired length. The prior art provides digging into the belt ends with barbs which can cause damage to the belt ends when the belt is in use. In addition, the prior art discloses a hollow belt which may undergo an excessive amount of flexting, compared with a solid belt, and the excessive flexing may contribute to excessive heat build up and premature wear. The prior art provides a connection between two ends of a belt in which the strength of the connection between the two ends of the belt may be overcome by the stresses of an operating engine, especially when the engine is started from a non-running condition. The foregoing disadvantages are overcome by the unique emergency drive belt apparatus of the present invention as will be made apparent from the following description thereof. Other advantages of the present invention over the prior art also will be rendered evident."}
{"text": "Field of the Invention\nThe present invention relates to a dynamics calculation method, a program and a recording medium used for a mechanism analysis simulation or the like. The present invention relates, more particularly, to a dynamics calculation method and a non-transitory recording medium that records a program suitable for use in dynamics calculation of a body made up of components connected through a movable axis and components connected through a fixed axis.\nDescription of the Related Art\nConventionally, when designing a machine product provided with an operation mechanism, mechanism analysis simulations are performed using a calculating machine such as a computer for confirming operation of the product in advance. In a mechanism analysis simulation, a movement of a whole product is simulated by performing dynamics calculation according to forward dynamics based on a dynamics model (also referred to as “simulation model”). A designer of the product (hereinafter referred to as “user”) can analyze not only operation of the product but also interference between components included in the product and a force applied to each component by performing a mechanism analysis simulation.\nIn the case of a multi-axis robot constructed by connecting many components through a movable axis or fixed axis, displacement/speed/acceleration of each component varies depending on, for example, torque input to each axis, and attitude in a time sequence is determined. However, since an actual movement of the multi-axis robot is affected by gravity/inertia or the like, it is difficult for a user to design a product by predicting such an affected movement. Thus, a link mechanism analysis apparatus and a link mechanism indirect data calculation apparatus are being proposed which perform simulations to confirm operation of a multi-axis robot having a link mechanism and confirm interference with a peripheral environment (Japanese Patent Publication No. 3361007). This apparatus simulates an overall movement by performing dynamics calculation using an amount of drive or drive force as an input parameter regarding a link mechanism.\nWhen performing a simulation using a computer, a user needs to construct a simulation model in advance by connecting many components through a movable axis or a fixed axis on a computer. A simulation model having a high degree of reusability is used when performing a simulation of a product such as a multi-axis robot system in particular in which replacement of components, change of component positions, change of a combination of components or the like frequently takes place. This model is convenient because components making up the multi-axis robot system are subdivided in advance and the user can thereby easily construct a simulation model in many patterns by only changing a combination of components as required.\nHowever, conventionally, the more the components are subdivided to increase the degree of reusability, the more noticeably the convenience as a simulation deteriorates. That is, when the components are subdivided, the number of components combined to construct a simulation model increases. Dynamics calculation has a feature that the amount of calculation increases according to the number of components. For this reason, when a simulation model formulated by combining many components is used, the processing load placed on dynamics calculation increases, consequently increasing a time required from the start of a simulation to the result output, resulting in poor usability.\nThe present invention has been implemented in view of the above-described problems and it is an object of the present invention to provide a dynamics calculation method, a program and a recording medium capable of performing high accuracy dynamics calculation without increasing a processing load when performing dynamics calculation based on a dynamics model formulated by combining many components."}
{"text": "The approaches described in this section are approaches that could be pursued, but not necessarily approaches that have been previously conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section.\nThe internet standard denoted MPLS (multi-protocol label switching) provides a data-carrying mechanism in which data packets are transmitted between network infrastructure elements, such as routers and switchers, according to labels identifying nodes and/or networks. A label is usually carried in a label stack entry of a data packet, and the entry usually comprises a 20-bit-long label field for storing the label.\nBecause the size of the label field in the data packet is limited, instead of encapsulating an actual Internet Protocol (IP) address of a node in the data packet, MPLS usually uses nicknames associated with the IP addresses. The nicknames may be distributed between routers using the Label Distribution Protocol (LDP), and the routers may translate the nicknames into the actual IP addresses or prefixes using a mapping. Unfortunately, retrieving the mapping and performing the translation may be time consuming and may impede network convergence time. Further, it may increase the complexity of LDP sessions, especially when an MPLS label needs to be exchanged between non-adjacent peers due to a link or node failure. Thus, using the nicknames as labels in MPLS, instead of actual addresses, may negatively impact network performance and convergence."}
{"text": "Autonomous machines and devices, such as autonomous robots, have been designed for performing various industrial and domestic functions. These domestic functions include lawn mowing, vacuum cleaning, floor sweeping and maintenance. By extending robots to these domestic functions, the person or user employing these robots has increased free or leisure time, as they do not have to expend the time required to perform the aforementioned tasks manually.\nThese autonomous robots typically operate in accordance with various computer programs that are part of the operating systems. Additionally, many of these autonomous robots are battery powered, and need to be charged once they are out of battery power. Additionally, if out of battery power, these autonomous robots typically stop where the power ran out and may be troublesome to locate or in difficult places to reach.\nAs a result, the autonomous robot must be located and manually brought to the charging unit, typically an electrical outlet. These processes require the user taking the time to perform them. Additional time is wasted as the user typically must wait a few hours before the robot is recharged, so it can start fresh again with fully charged batteries."}
{"text": "Asphalt-based roofing systems and products are well known. They include, for example, asphalt shingles and asphalt roll roofing. Many conventional materials are utilized as raw materials in the manufacture of asphalt roofing systems and products.\nAsphalt roofing systems and products generally comprise a substrate which is filled and coated with various asphalt materials. Generally, the substrate is filled with a \"saturant\" asphalt. A saturant asphalt is oil-rich and relatively nonviscous, to provide maximum waterproofing and saturation of the substrate. The saturant asphalt serves as a preservative, a waterproofing agent and an adhesive agent.\nThe saturated substrate is sealed by application of a harder, more viscous \"coating\" asphalt to both sides of the substrate. Coating asphalts generally contain finely divided minerals therein as stabilizers or fillers. Such compounds as silica, slate dust, talc, micaceous materials and dolomite have been utilized as fillers to render the coating asphalt more shatter-proof and shock-proof in cold weather.\nThe exterior, outer, or exposed surface of asphalt roofing systems and products is generally provided with a covering of granular material or roofing granules embedded within the coating asphalt. The granular material generally protects the underlying asphalt coating from damage due to exposure to light, in particular ultraviolet (UV) light. That is, the granules reflect light and protect the asphalt from deterioration by photodegradation. In addition, such granular material improves fire resistance and weathering characteristics. Further, colors or mixtures of colors of granular material may be selected for aesthetics.\nIn general, the mineral materials, particles or granules are embedded within the coating asphalt under pressure and are retained therein by adherence to the asphalt. With respect to each granule, the asphalt may be viewed as a \"hot sticky mud\" into which the granules are pressed. When the asphalt cools, pockets having the granules retained therein are formed.\nGood adherence of the roofing granules to the roofing product is beneficial. Loss of granules reduces the life of the roof, since it is associated with acceleration of photodegradation of the asphalt. In addition, the aesthetics of the roofing system may be compromised if granules are lost. Further, reduction of granule loss during installation improves safety conditions on the roof.\nGranule loss can also occur due to physical abrasion of the granular surface. This may occur any time a person walks on an installed roof for maintenance, during installation of the roofing surface or by such environmental conditions as tree branches rubbing on the granular surface and the physical contact of rain or hail with the roofing surface.\nIt has been found that adherence between the roofing granules and the coating asphalt is subject to deterioration by moisture. Granule-asphalt adhesion is not well understood. However, it is probable that secondary bonding interactions contribute to adhesive bond strength. Disruption of this secondary bonding by moisture may lead to decreased adhesion of granules to asphalt. Although water run-off from a slanted roof is generally sufficient to avoid prolonged exposure to moisture and thus to avoid substantial degradation by moisture to the granule/asphalt bond or interface, problems from moisture deterioration nevertheless pose substantial risk. For example, deterioration may be substantial in humid environments or in relatively flat portions of roofs where water can collect. Further, in many instances bundles of shingles (or similar roofing material) are stored in plastic wraps or containers prior to installation. Moisture trapped within such wraps or containers may cause substantial deterioration of the granule/asphalt bond, with resultant reduction in the integrity of the later installed roofing surface.\nPrior to applicants' improvements to the adhesion of roofing granules to the roofing product, it was generally felt that granule asphalt adhesion was satisfactory. It is, however, clear from the above discussion that beneficial results may be achieved by improving the granule asphalt adhesion in roofing products. What has been needed has been a method of improving asphalt-based roofing systems having granular material embedded therein with respect to granule loss due to moisture attack compromising the granule/asphalt bond or interface. In addition, improved roofing materials with respect to photodegradation of the asphalt layer by preventing granule loss by physical abrasion have been desired."}
{"text": "The laws regarding firearm safety and regulations are in a constant state of flux. In some cases a firearm that used to be legal may become illegal and to make it legal the owner must add certain features to the firearm to make it comply with the new regulations. In some cases the modifications can be performed with the use of a kit that can be used to modify the firearm."}
{"text": "The use of solid state drives (SSDs) in data systems is increasing. Solid state drives are advantageous over hard disk drives (HDD) because SSDs have no moving parts. Thus, SSDs are not as fragile as HDDs. Further, because SSDs have no mechanical delay, SSDs are not subject to the relatively lengthy access times, seek times and latency to which HDDs are subject. For this reason, SSDs are usually much faster than HDDs. A significant drawback of SSDs, however, is their relatively high cost compared to HDDs.\nSSDs can be based on either volatile (RAM-based) or non-volatile (flash) memory. Non-volatile memory can maintain information even when there is no power, whereas volatile memory requires continual power to maintain stored information. Flash memory stores information in memory cells that are made from floating-gate transistors. Solid state devices include single-level cell devices (SLC) and multi-level cell devices (MLC), for example. SLC devices store only one bit of information per cell, whereas MLC devices can store more than one bit of information per cell. Memory cells can be either NOR gate flash or NAND gate flash. NOR allows for random-access reading and NAND allows only page access.\nThere is a need for an improved data storage system and method that includes SSDs and utilizes data progression to improve the efficiency of the system."}
{"text": "Wireless communication devices with global location navigation system (GNSS) capability are becoming more prevalent. These devices depend on RF signals received from satellites for calculating position. However, these satellite signals are weak and are attenuated when inside buildings such that wireless devices can no longer obtain a lock on the signals and thus can no longer determine location.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "1. Technical Field\nThis disclosure relates generally to mobile device operations, and in particular, to handling data storage operations involving a mobile device.\n2. Description of the Related Art\nMobile devices (e.g., personal digital assistants (PDAs), cellular phones, mobile messaging devices, etc.) frequently perform data storage operations in order to store data used on a mobile device. The data to be stored can be information from many different types of sources, such as from the user of the mobile device, software applications operating on or externally to the mobile device. Difficulties can arise in attempting to handle data storage operations in an efficient and robust manner. The difficulties are enhanced on a resource constrained device, such as a mobile device, where a traditional file system may be deemed as being too “heavyweight” for such a device."}
{"text": "It is a recognized fact that electronic devices, and especially BJT transistors, should be operated within safe limits for proper operation and without the risk of becoming damaged.\nIn many applications, such devices are operated far from the limiting conditions; however, for some of them, such as power devices, it is extremely important to bring component performances to the limit, but without overtaking the limiting conditions which would result in the device being damaged.\nShown on a graph in FIG. 1 are the working points in safe limit conditions of a power transistor (FBSOA (Forward Base Safe Operating Area) curve). In particular, on the graph, a set of collector current curves are plotted against the collector-emitter voltage for a same power device as the duration of the base current varies.\nAs follows from this FIG. 1, if the base current is of the pulsed type, as the duration of the individual pulses decreases (100 msec, 1 msec, 100 .mu.sec, 10 m sec, 1 .mu.sec), the area included by the FBSOA curve (curves I, II, III, IV, V) increases.\nIn order to limit the maximum dissipatable power during operation, protection circuits are often associated with such devices. FIG. 1 also shows a curve (VI) of possible operation of a linear protection circuit according to the prior art.\nA first embodiment of that protection circuit is shown in FIG. 2, which depicts an electronic regulator circuit for driving a power device Q0. That regulator circuit comprises a first driving portion 1 and a second protection portion 2.\nThat first driving portion 1 is realized by a first amplifier OP1 with a first input terminal I1 connected to a voltage reference Vref, a second input terminal I2 connected to a voltage divider Vd, and an output terminal O1 connected to a control terminal 6 of the device Q0.\nThe second protection portion 2 comprises a second amplifier OP2 having a feedback resistor Rs connected between its input terminals I3 and I4.\nIn particular, the amplifier OP2 is an operational stage with an inherent offset voltage Voffset.\nThe input I3 is connected to a terminal 4 of the power device Q0 through a series connection of a Zener diode Dz and a resistor R. An output terminal O2 of the second amplifier OP2 is connected to the terminal 6 of the device Q0.\nOne end of the resistor Rs is also connected to a terminal 5 of the device Q0, the other end being connected to the divider Vd.\nIn the circuit architecture of FIG. 2, the current flowing in the output load goes through the resistor Rs and is picked up at the input terminals of the amplifier OP2.\nSo long as the voltage picked up across the resistor Rs is lower than the voltage Voffset of the amplifier OP2, the output voltage of the power device Q0 will be regulated through the driving portion 1.\nUpon the load current exceeding the maximum current for which the transistor Q0 has been designed, the increased voltage Voffset will set the protection circuit to operate.\nA more detailed description of the operation of a regulator circuit of that type is found in U.S. Pat. No. 5,714,905, which is incorporated herein by reference.\nWhile being in many ways advantageous, this solution has several drawbacks.\nAs the fabrication process varies, the operation curve of the protection circuit SOA can vary to the point of overtaking the FBSOA curve. To prevent the power device from operating outside the safe operation range, the power device is provided oversize such that the SOA curve is contained within the FBSOA curve with ample margin.\nA first disadvantage of this solution is that the capacity for integration of the power device is altered for the worse, because its being oversize involves the use of a larger amount of silicon for its formation. This obviously results in an undesired cost increase.\nA second disadvantage of this construction is the appearance of the so-called latch-down phenomenon in the regulator circuit. With high values of the power device working voltage, the protection circuit heavily limits the current delivered from Q0, which may create difficulties in initially charging the capacitive loads provided downstream of the circuit."}
{"text": "The invention is directed to new plastisols with improved qualities, as well as to the use of advantageous (meth)acrylate resin masses for the production of these new plastisols.\nPlastisols are generally found as a two-phase system having one component as a plastic (binding material) and the second as a suitable softening agent. As a matter of principle, a large variety of different plastics are usable as binding materials. However, from a technical point of view, only rather few plastics are actually utilized. By far, the most important class of polymers that is utilized for this purpose are derived from polyvinyl chloride (PVC). However, from the point of view of environment considerations, the utilization of PVC is undesirable. The danger of dioxin formation in the case of fire and the corresponding contamination of the surroundings militates against the utilization of PVC.\nFor this reason, attempts have been made to create plastisols based on poly(meth)acrylates. See for example, DE-PS 934 498, FR-A 2,291,248 and EP 0 774 483 A2. As used herein, the terminology (meth)acrylates includes both acrylates and methacrylates.\nA fundamental problem associated with the use of poly(meth)acrylate plastisols in comparison to PVC plastisols rests in their insufficient shelf-life and mechanical properties. The use of spray-dried emulsion polymerisates based on PMMA in accordance with the state of the art, leads in combination with commercially available softening agents to good gelation and film properties, but also to reduced storage stabilityxe2x80x94the products already gel after comparatively short time at room temperature which can be recognized by an increase in viscosity.\nThe problem of storage capacity was attempted to be solved as disclosed in EP 0 774 483, by raising the amount of the middle sized particles by providing milled-suspension polymerisates to the usual emulsion polymerisates. The disadvantage of this procedure is in the circumstance that the higher particle size leads to impeded gelling of the plastisols, since under certain circumstances, the larger suspension polymerisates no longer gel completely which this technology demands. The inhomogeneity of the films caused by the incomplete gelation gives rise to insufficient mechanical properties with respect to tensile strength (resistance to tearing) in the corresponding films. Furthermore, one must not ignore the optical disadvantage of these inhomogeneous films in the formation of smooth surfaces. In order to achieve the production and reduction in size of the suspension polymerisates, additional procedural steps are required which lead to increase in cost of the product.\nFurthermore, the art similarly requires the lowest possible viscosity level for the plastisol. On the one hand, it is desirable to cover the broadest possible spectrum of processing techniques and thus cover the widest possible area of usage. On the other hand, plastisols should desirably be processable at the lowest possible temperature and only after warming and subsequently cooling, should the plastisol gel to a solid uniform film.\nThe state of the art describes the use of emulsion polymerisates up to a molecular weight of 2,000,000 g/mol (EP 0 539 031 A1). For technically usual PMMA plastisols, there are utilized emulsion polymerisates of a substantially lower (less than 500,000 g/mol) molecular weight.\nHeretofore, the state of the art does not describe polymerisates and/or copolymerisates of the utilized (meth)acrylate with such high chain lengths for the formation of plastisols. Thus, there is no suggestion to one skilled in the art that by utilizing polymerisates and/or copolymerisates of (meth)acrylates of such high molecular weight, one could obtain the above-described advantageous plastisols. On the contrary, one skilled in the art would be under the impression that by raising the molecular weight and therewith chain length polymers, there would be obtained an increase in the viscosity of the plastisols. In fact, as is shown in Table 1 hereof, what occurs is the exact reverse.\nThe task of the present invention is to find a plastisol that combines good film qualities and gel formability which lead to good tensile strength and break/extension properties of the product with acceptable storage stability and low viscosity. These and further named tasks, which will be clear to one skilled in the art from the state of the art, are solved by plastisol defined by claim 1 hereof.\nAdvantageous embodiments of the invention are set forth in claims relating back to claim 1. Uses of plastisols in accordance with the present invention are described in the corresponding claims.\nThese plastisols show:\n1. Polymerisates and/or copolymerisates of (meth)acrylates obtainable by the polymerization of mixtures containing, as polymerizable components:\n(A) 20 to 100 wt. % methylmethacrylate.\n(B) 0 to 80 wt. % of (meth)acrylates different from methylmethacrylate, having formula I: \nxe2x80x83wherein\nR1 is hydrogen or methyl, and\nR2 is a linear or branched (C1-C18) alkyl residue.\n(C) 0 to 40 wt. % of further monomers different from (A) and (B),\n(D) 0 to 40 wt. % of adhesion causing monomer, wherein (A) through (D) yield 100 wt. % of the polymerizable components.\nII. A compatible softening agent in the proportion of 5 to 400 parts by weight, relative to 100 wt. % of polymerisate and/or copolymerisate I.\nIII. Inorganic fillers in the amount of 9 to 700 parts by weight, relative to 100 parts by weight of component I, characterized thereby that the mean molecular weight Mw of the poly-merisate and/or copolymerisate of the (meth)acrylate is greater than 3,500,000 g/mol, enables one to provide, a particularly advantageous and yet surprisingly stable plastisol having low viscosity which, because of homogeneous film formation on gelation, show excellent mechanical as well as optical properties. The polymerisate and/or copolymerisates are preferably obtained through emulsion polymerization.\nUnder normal circumstances, a rise in molecular weight, that is to say chain length, in solutions or emulsions of polymer, leads to a faster separation of mixtures and thus a poorer storage life. However, contrary to expectations, the storage stability of plastisols improves with the increase in chain length. In the plastisols of the present invention, one observes a far lower level of mixture separation than with plastisols formed with polymerisates of the state of the art.\nEspecially preferred in the plastisols of the present invention based on (meth)acrylate, are those wherein the residue R2 of the (meth)acrylate of formula I is a linear or branch chained (C1-C8) alkyl residue.\nIt is particularly advantageous if the mean molecular weight Mw of polymerisate and/or copolymerisate of the charged (meth)acrylate in the plastisols is greater than 3,500,000 g/mol., suitably greater than 3,900,000 g/mol.\nParticularly desirable for the formation of plastisols is the use of polymerisate and/or copolymerisate of the (meth)acrylate with a mean molecular weight Mw of greater than 3,500,000 g/mol, preferably greater than 3,900,000 g/mol.\nThe term xe2x80x9c(meth)acrylatexe2x80x9d within the scope of the present invention includes both acrylates and methacrylates.\nThe chain length of the polymerisates and/or copolymerisates of the methacrylates are limited by a synthetically achievable chain lengths. Chain lengths having a mean molecular weight of Mw of about 12,000,000 g/mol are achievable.\nThe mean molecular weight Mw of polymers for use of the present invention are determined by means of SEC or GPC (size exclusion chromatography or gel permeation chromatography) relative to polystyrene standards. Those skilled in the art are familiar with SEC or GPC analysis methods for the determination of the mean molecular weight.\nWithin the scope of the invention, a further quantity for the characterization of the molecular mass of the charge polymerisate and/or copolymerisate is the viscosity number VZ. This viscosity number is determined by the standard set forth in DIN 51 562 Parts 1 and 3 (Revision of January 1983 and May 1985).\nThe polymerisates and/or copolymerisates used in the plastisols of the present invention have a VZ value equal to, or greater than 600, desirably greater than 600, preferably greater than 900, and especially preferably greater than 1200.\nBy the term xe2x80x9clinearxe2x80x9d or xe2x80x9cbranch chain (C1-C18) alkyl residuexe2x80x9d, one understands a group of alkyl residues commencing with methyl via ethyl, up to an 18 carbon atoms containing radical. Included in this group are several bonding isomers which are conceivable within the group.\nIncluded in the different monomers of (A) and (B) are styrene and its derivatives, vinyl esters such as vinyl acetate, vinyl propionate, vinyl esters of higher alkyl acids, maleic acid anhydride, itaconic acid and its esters, and olefins such as ethylene, propylene, isobutylene, etc.\nAs adhesion improving monomers (HM) may be included radically polymerizable monomers which contain functional groups, which can enter into interactions with materials that are to be coated. Such interactions can occur by, for example, hydrogen bonding, complexing, dipole forces and the rest, in which, generally speaking, hetero atoms such as oxygen or nitrogen take part. There may be mentioned as functional groups amino, in particular dialkyl amino, (cyclic)amido- imido- hydroxy- oxy- carboxyl- and cyano groups. Such monomers are, generally speaking, known (see H. Rauch Puntigam, Th. Volker, (translated) Acryl- and Methacryl-compounds, Springer-Verlag, 1967; Kirk-Othmer, Encyclopedia of Chemical Technology, 3rd. Ed., Vol. 1, pp. 394-400, J. Wiley 1978; DE-A 25 56 080; DE-A 26 34 003).\nPreferably, the adhesion improving monomers HM belong to a nitrogen containing class of vinyl heterocycles, suitably having 5 or 6 rings and/or copolymerizable vinyl carboxylic acids and/or hydroxyalkyl-, alkoxy-alkyl-, and aminoalkyl-substituted esters or amides of acrylic and methacrylic acids.\nAs nitrogen containing heterocyclic monomers of the HM group, particularly desirable are the class of vinyl imidozoles, vinyl lactams, vinyl carbazoles and vinyl pyridines. As examples, which no way should be considered to be limiting for these monomeric imidazole compounds may be mentioned N-vinyl-imidazole, (also called vinyl-1-imidazole), N-vinyl-methyl-2-imidazole. N-vinyl-ethyl-2-imidazole, N-vinyl-phenyl-2-imidazole, N-vinyl-dimethyl-2,4-imidazole, N-vinyl-benzimidazole, N-vinylimidazoline (also called vinyl-1-imidazoline), N-vinyl-methyl-2-imidazoline, N-vinyl-phenyl-2-imidazoline and vinyl-2-imidazole.\nAs examples of monomers which are obtained from lactams, there may be mentioned in particular N-Vinylpyrrolidone, N-Vinylmethyl-5-pyrrolidone, N-Vinylmethyl-3-pyrrolidone, N-Vinylethyl-5-pyrrolidone, N-Vinyldimethyl-5,5-pyrrolidone, N-Vinylphenyl-5-pyrrolidone, N-Allyl-pyrrolidone, N-Vinylthiopyrrolidone, N-Vinylpiperidone, N-Vinyldiethyl-6,6-piperidone, N-Vinylcaprolactam, N-Vinylmethyl-7-caprolactam, N-Vinyl-ethyl-7-caprolactam, N-Vinyl-dimethyl-7,7-caprolactam, N-Allylcaprolactame, N-Vinylcapryllactam.\nUnder monomers which may be derived from carbazole, there may be particular mentioned N-Vinylcarbazole, N-Allylcarbazole, N-Butenylcarbazole, N-Hexenylcarbazol and N-(Methyl-1-ethylene)carbazole. As copolymerizable vinylcarboxylic acids, there may be mentioned as particularly suitable acrylic and methacrylic acids as well as suitable salts thereof.\nFurther, there may be utilized as the following oxy, as well as alkoxy substituted alkyl esters of (meth)acrylic acids, in particular 2-Hydroxyethyl (meth)acrylate, Hydroxypropyl(meth)acrylate, 2-Methoxy-ethyl(meth)acrylate, 2-Ethoxyethyl(meth)acrylate, 2-Butoxyethyl(meth)acrylate, 2-(2-Butoxyethoxy)ethylmethacrylate, 2-(Ethoxyethyloxy)ethyl(meth)acrylate, 4-Hydroxybutyl(meth)acrylate, 2-[2-(2-Ethoxy-ethoxy)ethoxy]-ethyl(meth)acrylate, 3-Methoxybutyl-1-(meth)acrylate, 2-Alkoxymethylethyl(meth)acrylate, and 2-Hexoxyethyl(meth)acrylate.\nFurthermore, the following amino substituted alkyl esters (meth)acrylic acids, may be mentioned 2-Dimethylamino-ethyl (meth)acrylate, 2-Diethylaminoethyl-(meth)acrylate, 3-Dimethylamino-2,2-dimethylpropyl-1-(meth)acrylate, 3-Dimethylamino-2,2-dimethylpropyl-1-(meth)acrylate, 2-Morpholinoethyl(meth)-acrylate, 2-tert.-Butylamino-ethyl(meth)acrylate, 3-(Dimethylamino)propyl-(meth)acrylate, 2-(Dimethyl-aminoethoxyethyl)-(meth)acrylate.\nAs examples of (meth)acrylamide, the following monomers may be mentioned N-Methyl(meth)acrylamide, N-Dimethylaminoethyl(meth)-acryl amide, N-Dimethylaminopropyl(meth)acrylamide, N-Isopropyl(meth)-acryl amide, N-tert.-Butyl(meth)acrylamide, N-Isobutyl(meth)acrylamide, N-Decyl(meth)acrylamide, N-Cyclohexyl(meth)acrylamide, N-[3-(Dimethyl-amino)2,2-dimethylpropyl]-methacrylamide, N-[2-Hydroxyethyl](meth)-acrylamide.\nThe plastisols of this application contain generally known softeners W in the usual amounts, for example for PMMA plastisols (but also for PVC plastisols), in particular phthalic acid esters, adipic acid and/or sebacic acid esters, chlorinated paraffins, trialkyl phosphates, aliphatic or araliphatic polyesters, as well as various other polymeric softening agents such as urea soft resin (see H. K. Felger (Trans.), Plastics-Handbook Vol. 1/1C, Kanset-Verlag 1985 as well as in H. F. Mark, et al., Encyclopedia of Polymer Science and Engineering, Supplemental Volume, pg. 568-647, J. Wiley 1989). A selection of suitable softening agents can also be found in DE-C 25 43 542. Similarly, there may be used softening agents suitably softening agent combinations which, during a storage of the prepared plastisols for 3 weeks at 30xc2x0 C., limit a rise in viscosity to 10 times, suitably only 5 times the initial viscosity. Particularly may be mentioned dioctylphthalate, diisodecyl-phthalate, diethylhexylphthalate, Di-C7C11-n-alkylphathlate, tricresylphos-phate, dibenzyltoluene (LIPINOL?T, Product of Fa. Huls AG), 2,2,4-trimethyl-1,3-pentandioldibenzoate (Benzoflex 354, Product of Velsicol) and benzyl-octytphtalate.\nFurthermore, the plastisols may contain generally speaking, well known inorganic fillers in amounts in amounts up to 700 parts by weight (per 100 parts by weight of plastisol). There may be mentioned for example, calcium carbonate (chalk) titanium dioxide, calcium oxide, Perlite, precipitated and coated chalks as rheologically active additives, further similarly thixotroping agents such as, for example, pyrogenic silica acid. The particle size generally lies in the range of 5 to 25 m and furthermore the plastisols according to the intended use, may contain known active agents such as adhesive agents, cross-linking agents, stabilizers, lubricating agents, blowing agents in amounts of 0-5 wt. % (relative to the plastisols). There may also be utilized calcium stearate as a lubricating agent.\nThe plastisols of the present invention may be utilized for a variety of purposes such as those designated for PVC or PMMA plastisols. In particular, there may be considered walking surfaces for floor coverings, free transparent folios, KFZ underlayment protection, corrosion protection for metals, intermediate coverings for floor coverings (foamed), wallpaper, plastic sheets, car roof windows, instrument panel, coverings, crown corks.\nBasically, the above-identified monomers, in particular methylmeth-acrylate as well as the monomers of Group B are useful for the core-shell polymerisates.\nThe construction of the polymerisate and/or copolymerisate from a core material and a shell material proceeds in a generally known manner in a particular procedure by emulsion polymerization. In this process, the monomers forming the core material are polymerized in aqueous emulsion in a first procedural step. When the monomers of the first step are substantially fully polymerized, the monomeric components of the shell material are added under such conditions to the emulsion polymerisate that formation of new particles is avoided. Thus, during the second step the thus formed polymerisate forms a shell around the core material.\nIn the first polymerization step in which the core material is formed, between 0.01 to 3 wt. % (relative to the monomer) of an anionic cationic or nonionic emulsifier such as sodium lauryl sulfate, alkyl benzyl sulfonate, oxy ethylation products of alkylated phenols or their neutralized sulfonation products are utilized. The polymerization suitably proceeds between 60xc2x0 and 100xc2x0 C. with the help of a water soluble radical source such as potassium or ammonium persulfate or hydrogen peroxide. Before the commencement of the second polymerization step, further initiator may be added. However, very mall amounts or in fact no additional emulsifier is used in the second step."}
{"text": "Trocars are generally used as access devices for introduction of other surgical instruments into a patient. A conventional thoracic trocar is disclosed in U.S. Pat. No. 5,279,575 to Sugarbaker. In U.S. Pat. No. 5,279,575 the trocar has a pair of flukes which are pivotally coupled to the body with pins. The trocar also includes a rigid clamp which is slidably on the outer tube for securing the trocar to a body wall. The inner tube slides longitudinally relative to the outer tube for moving the flukes to the open position. A distal end of the inner tube engages a camming surface on each of the flukes for camming the flukes outward.\nThe trocar disclosed in U.S. Pat. No. 5,279,575 suffers from a number of drawbacks. First, the clamp is relatively stiff and does not conform to the shape of the patient. As such, the trocar may be difficult to tilt at a desired angle when passing between adjacent ribs. Another drawback of the trocar is that the hinges create a relatively large profile since the hinges extend beyond the outer, cylindrical surface of the tube. The relatively large profile increases trauma to the patient which is particularly problematic when the trocar is rotated or otherwise displaced after initial placement. Yet another drawback of U.S. Pat. No. 5,279,575 is that the clamp is mounted to a third tube which slides over the outer tube. The three tube configuration also increases the profile of the device and trauma to the patient."}
{"text": "The invention relates in general to a medical device structured to serve as a replacement for damaged, diseased, pathologic, degenerated or ruptured intervertebral discs, and more particularly to such a device which is self-centering and maintains the proper spatial relationship between adjacent vertebrae while allowing for movement of the adjacent vertebrae in a manner similar to and consistent with the range and type of motion found prior to pathology, degeneration or rupture.\nIn addressing the problems of pathologic, degenerated or ruptured intervertebral discs, the solutions have generally focused on immobilizing the adjacent vertebrae by fusion or replacement of the damaged disc with a prosthetic disc, where the prothesis is designed to allow for some motion between the adjacent vertebrae. A successful disc replacement requires a prosthetic device which remains properly positioned between the vertebrae, maintains the vertebrae in proper spaced relation, allows for a range of motion while maintaining proper alignment of the spine and vertebrae in the frontal and sagittal plane during such motion, is biocompatible and is of sufficient durability to withstand loading cycles of up to 100 million or more.\nThere have been numerous patents issued in this field. Some patents involve replacement discs composed of a compressible or elastic material, some involve discs comprising matrices to enhance bone growth into the discs, some involve discs composed of both rigid and non-rigid materials, and some involve rigid discs. Of the later type, some are designed to remain detached from the adjacent vertebrae while others are designed for direct connection to the vertebrae. Multi-component replacement discs are known where the plural components are free to move relative to each other. The major drawback from the known devices is that the range, angle and type of movement does not properly mimic the natural range, angle and type of motion found in the actual human disc, while maintaining proper frontal and sagittal alignment of the spine.\nIt is an object of the replacement intervertebral disc of the present invention to provide a device which substitutes for a human disc and successfully replicates the range, angle and type of movement between the adjacent vertebrae, where the disc is composed of a durable, rigid, biocompatible material of high fatigue strength so as to withstand millions of loading cycles without failure or degradation. It is a further object to provide such a disc having two generally convex surfaces which generally correspond to the disc being replaced and the concave shape of adjacent vertebrae, where the maximum dimension between the two convex surfaces is not central to the main body of the disc, where the disc is self-centering to preclude disc impingement into the spinal canal, and where proper intervertebral spacing is maintained by joining the two opposing convex surfaces by a generally vertically oriented wall member. It is a further object to provide such a disc where the convex surfaces are not in parallel relationship so as to better replicate the configuration of a human disc and proper spinal alignment."}
{"text": "1. Field of the Invention\nThe present invention relates to small, chain-driven carts (sometimes referred to as \"go-carts\"), and the like having an internal combustion engine of, for example, five-ten (5-10) horse power that drives the cart rear wheels via an axle and centrifugal clutch. Even more particularly, the present invention relates to an improved engine shaft, clutch, and chain-driven axle assembly for powering small motorized carts wherein an improved lubrication arrangement allows grease to be injected through an engine drive shaft (or a short intermediate shaft or \"jack\" shaft) of the cart via a porting system so that the clutch is lubricated when rotating at high speeds, namely, when the cart is at rest and the engine idles.\n2. General Background\nSmall chain-driven motorized carts are commercially available from a number of manufacturers. These carts are used primarily for off road motoring at relatively low speeds. They are frequently used for backyard and off road family fun, and can be found at amusement parks and fairs on occasion.\nThese small chain-driven carts include a metal frame that is generally of welded steel, using square and round tubing and plate steel, for example. Such small carts use an internal combustion engine for power, usually on the order of five to ten (5-10) horse power, as an example. The cart frame is constructed very low to the ground for stability purposes. The engine on such a cart is usually a commercially available internal combustion engine, having a drive shaft that carries a centrifugal clutch. A sprocket on the clutch engages an endless belt or chain. Another sprocket rigidly affixed to the axle also engages the endless belt or chain. When the engine of the cart idles, the centrifugal clutch by design is disengaged. When the user presses a throttle (e.g., foot pedal) the engine revolutions increase and the clutch engages so that the sprocket of the clutch rotates and drives the endless belt or chain and its associated axle sprocket and attached axle.\nBy placing the cart low to the ground, equipping the cart with relatively small wheels, and by using the cart in a off road environment, the centrifugal clutch is subjected to dirt, sand, dust, and debris which accumulates in and around the clutch assembly. This produces an abrasive environment which can quickly wear out the clutch parts. The centrifugal clutch is typically rotated at high speeds when the cart idles. The clutch engages upon increasing engine revolutions so that centrifugal members engage the clutch and it forms a rigid connection with the drive shaft of the cart.\nTypically, these carts use a horizontal shaft engine. The drive shaft extends horizontally from one side of the engine. However, it should be understood that a short separate intermediate drive shaft sometimes referred to a \"jack-shaft\" can be used to drive the carts. A \"jack-shaft\" is simply a short shaft mounted between bearings and in spaced-apart position away from the horizontal shaft of the engine. The \"jack-shaft\" is driven by the motor shaft using a belt or sprockets and a chain, for example. In either situation, the present invention is designed to provide an improved lubrication arrangement for the engine shaft, and/or a \"jack-shaft\" used with the engine drive."}
{"text": "1. Field of the Invention\nThis invention relates to microelectromechanical systems for use as interferometric modulators. More particularly, this invention relates to improved methods of manufacturing microelectromechanical system devices having different sized cavities between a movable element and a substrate.\n2. Description of the Related Art\nMicroelectromechanical systems (MEMS) include micro mechanical elements, actuators, and electronics. Micromechanical elements may be created using deposition, etching, and or other micromachining processes that etch away parts of substrates and/or deposited material layers or that add layers to form electrical and electromechanical devices. One type of MEMS device is called an interferometric modulator. As used herein, the term interferometric modulator or interferometric light modulator refers to a device that selectively absorbs and/or reflects light using the principles of optical interference. In certain embodiments, an interferometric modulator may comprise a pair of conductive plates, one or both of which may be transparent and/or reflective in whole or part and capable of relative motion upon application of an appropriate electrical signal. In a particular embodiment, one plate may comprise a stationary layer deposited on a substrate and the other plate may comprise a metallic membrane separated from the stationary layer by an air gap. As described herein in more detail, the position of one plate in relation to another can change the optical interference of light incident on the interferometric modulator. Such devices have a wide range of applications, and it would be beneficial in the art to utilize and/or modify the characteristics of these types of devices so that their features can be exploited in improving existing products and creating new products that have not yet been developed."}
{"text": "Measuring instruments are frquently provided on the handlebar of a bicycle for indicating to the cyclist relevant operational data like the actual speed or the covered distance. These instruments are detachably mounted on the handlebar so that they can be removed particularly in view of the danger of theft when the bicycle is parked in a place accessible to the public. The electronic instrument is supplied with signals from a signal receiver which is mounted on a fork leg of a bicycle fork unit carrying a bicycle wheel. The signals are generated in the signal receiver by a signal inducer fastened to a spoke of the respective wheel. In order to permit easy detachment of the instrument from the instrument carrier fixed to the handlebar the cable is connected to the instrument carrier and is in electrical connection with electrical contacts provided on the instrument carrier. Counter-contacts are provided on the instrument. The counter-contacts are electrically connected to the contacts when the instrument is engaged with the instrument carrier."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical mount.\n2. Background Information\nMechanical mounts are commonly used to secure and vary the position of an optical element such as a mirror. The mount may be part of an optical system such as a laser or an interferometer. The system may be tuned or calibrated by varying the incident angle of the light reflected from the mirror. A conventional optical mount typically has a number of adjustment screws that can be manipulated to tilt the mirror and vary the angle of incidence.\nThere have been marketed gimbal optical mounts that can tilt the reflecting mirror and change the incident angle without spatially shifting the point at which the light is reflected from the mirror. This can be accomplished by reflecting light from the pivot points of the gimbal mount which are coincident with the mirror face. Gimbal mounts allow a scientist to vary the incident angle without introducing a phase shift in the reflected light.\nConventional gimbal optical mounts typically contain two concentric rings. The innermost ring holds the optical element and can rotate relative to an outer ring. The outer ring can rotate relative to an outer stationary housing about an axis that is perpendicular to the axis of rotation of the inner ring. Ring type gimbal mounts can be relatively expensive to produce.\nU.S. Pat. No. 4,088,396 issued Edelstein and U.S. Pat. No. 3,357,268 issued to Richter disclose optical gimbal mounts which each have a first plate and a second plate. The first plate is typically mounted to a stationary surface such as a table. The second plate holds the optical element and has a spherical bearing surface. The spherical bearing surface cooperates with a spherical seat of the first plate to allow the second plate and optical element to pivot about a point. Unfortunately, the pivot point is located away from the optical surface of the optical element so that pivoting the element may introduce a phase shift in the reflected light.\nThere have been published articles which discuss the use of gimbal optical mounts that have kinematic type interfaces which pivot the optical element about a point on the optical surface. These mounts have three or four adjustment screws that are manipulated to pivot the optical element. Accurately pivoting the optical element typically requires manipulating two screws the same amount. Manipulating two screws can be both cumbersome and time consuming. For example, a gimbal mount with three or four screws requires the simultaneous manipulation of at least two screws to pivot the optical element about an axis. It would be desirable to provide an opto-mechanical mount that can tilt an optical element without introducing a spatial phase shift and which can pivot the element about an axis by manipulating only one screw.\nOne embodiment of the present invention is an optical mount for an optical element which has an optical surface. The mount includes a first plate, and a second plate which holds the optical element. The second plate may have a spherical bearing surface that cooperates with the first plate so that the optical element pivots about a point on the optical surface. The mount has a single first adjustment device to pivot the optical element about one axis and a single second adjustment device to pivot the element about a second perpendicular axis."}
{"text": "1. Field of the Invention\nThe present invention is a method and apparatus for implementing smart management of a rechargeable battery. The invention uses the operation data of the rechargeable battery including the voltage, temperature and current flow to manage charging and discharging of the battery in an intelligent manner. The method is a dynamic technique for the management of the rechargeable battery that is able to make appropriate adjustments in response to any change in parameter values of the battery.\n2. Description of Related Art\nUse of high performance and long life batteries, low power consumption and product miniaturization are on-going objectives in the design of portable electronic devices such as notebook computers, personal digital assistants (PDA) and mobile phones. To further enhance the portability of an electronic product, the first thing a product designer must do is to find a long-life rechargeable battery that is capable of supporting the operation of the product for ten hours or more. However, the miniaturization of electronic products has severely constrained the storage capacity for a battery without sufficient space to accommodate additional rechargeable cells. In addition, designers are trying to package more functions in existing products to make them super versatile. Such devices with increased complexity after incorporating all the additional functions definitely need larger load current than their conventional counterparts. Developing a smart battery management system for rechargeable batteries to cut down unnecessary power loss has become imperative.\nPortable electronic products such as notebook computers usually have a built-in battery control circuit to optimize and adjust the current flow during the charging and discharging of the battery. In the conventional method, the control circuit in portable devices works in conjunction with a signal I/O interface to retrieve the parametric values from the battery. These periodically collected data relating to the battery operation are shown to the user of a portable device to make him or her aware of the current operational status of the battery.\nThe control circuit in conventional battery management is set up with predetermined values for the battery operation parameters. These parameters may be the lower or upper reference voltage, or lower or upper reference temperature. The control circuit then uses the data to control the charge and discharge of the battery. However, a rechargeable battery is usually made of chemical compounds such as a lithium-polymer or nickel-metal. During the operation of a battery, electrochemical reaction takes place simultaneously inside the battery. In the electrochemical process, the basic elements often exhibit non-linear characteristics. Simply by presetting the values of the external control factors such as operating temperature, flow of current during charge or discharge, and the time the battery being used is not adequate to control the operation characteristics of the battery with the design specs.\nSpecifically, changing in temperature affects the internal ohmic resistance of an electrode in the battery. Higher temperature could cause an increase of internal ohmic resistance, and vice versa, leading to unpredictable variation of the terminal voltage of the battery. Furthermore, longer battery use time could cause an increase in internal ohmic resistance, resulting in improper estimates of the upper and lower voltage limits.\nIn conventional methods, constant current flow is often used in the process of charging or discharging a battery. According to the charge/discharge characteristic curve of a battery, using a small current to charge a battery is more efficient than a high constant current. Also, using a large current to charge/discharge a battery could produce early saturation or depletion. This phenomenon can be explained by an increasing amount of electric energy being transformed to thermal energy as the battery approaches saturation. It is not economical to charge a battery with a constant current, since this electric energy cannot be efficiently saved in the storage battery. Besides, low ambient temperature or long idling of the battery could cause the electrolyte in the battery to coagulate or produce sedimentation. Under such conditions, charging with excessive current flow could easily damage the battery.\nIn the conventional method, if the battery is close to depletion in a discharge mode, only a message is issued by the control system to warn users of impending depletion, but the discharging process continues without interruption. As a result, the battery could be over discharged and create irreversible damage, such that the battery cannot be used anymore.\nThe main objective of the present invention is to provide a method and apparatus for implementing a smart battery management system on a rechargeable battery. The charging of a battery is carried out in two stages with the assistance of an optimizing charge controller. In the first stage, charging a battery is conducted in the normal mode with a constant flow of current. As the charging of the battery approaches saturation, the charge controller is switched to pulsation-charge mode to charge the battery intermittently. With the smart control, the charging efficiency of a rechargeable battery is improved notably.\nThe second objective of the present invention is to provide a trickle charger. If the ambient temperature of the battery is below a predetermined lower temperature limit or the terminal voltage is below a predetermined voltage lower limit, the trickle charger is enabled to slowly charge with a small current. The trickle charger serves to activate the electrolyte in the battery in preparation for a normal charge. These control methods can prevent possible harmful charging of the battery with excessive current.\nThe third objective of the present invention is to provide a discharge controller to protect the battery against over discharge. If the central processing unit detects that the voltage of the battery is below a predetermined lower voltage limit, representing excessive discharge from a battery, it automatically turns off the discharge controller to stop the discharging of the battery and cuts off all load supply current from the battery control circuit, entering a self-cutoff mode.\nThe fourth objective of the present invention is to provide a cutoff switch. If the central processing unit detects that the battery is not active, it cuts off the load current to the temperature detector to reduce power loss from the battery.\nThe apparatus necessary for implementing smart management of the rechargeable battery includes a central processing unit, a charge controller, trickle charger, a discharge controller, a cutoff switch and an economizer switch.\nThe central processing unit is respectively connected to a current flow detector for sensing charge/discharge current flow, a temperature detector for sensing the surface temperature of the battery, and a voltage detector for sensing the terminal voltage across the battery. These parametric values relating to the operational status of the battery collected by the above detectors are saved in an EEPROM.\nThe optimizing charge controller is connected across the central processing unit and the rechargeable battery. If the battery operates within the normal range of terminal voltage and temperature, the central processing unit activates the charge controller to carry out normal-charge of a battery with a constant current.\nThe trickle charger is connected across the central processing unit and rechargeable battery. In the charge mode, it can provide slow charge of the battery with a small current to activate the electrolyte in the battery before the normal charging process is initiated.\nThe discharge controller is connected across the central processing unit and the rechargeable battery. If the terminal voltage of the battery is below a predetermined lower voltage limit during discharging, the central processing unit hereupon issues a control signal to turn off the discharge controller stopping the discharging and turns off the voltage regulator entering a self-cutoff mode.\nThe cutoff switch is connected to the temperature detector and central processing unit. If the central processing unit detects that the battery is idling, the processor turns off the temperature detector to reduce power loss.\nThe economizer switch is connected to the rechargeable battery and the central processing unit through the voltage regulator. If the battery is connected to an external power source for recharging, the economizer switch activates the voltage regulator which in turn supplies normal operating voltage to force the system to return to normal operation from the previous self-cutoff mode.\nThe features and related techniques of the present invention will become more apparent from the description of the preferred embodiment when taken with the FIGS. of the invention."}
{"text": "The present invention relates to gloves and, more particularly to waterproof gloves having abrasive scouring surfaces for purposes of cleaning.\nHeretofore, waterproofed gloves and mittens, have been used to protect the user's hands from harsh detergents and other chemicals necessary in most cleaning operations. In addition, steel wool and other such fiberous abrasives have been used to accomplish the cleaning of surfaces of various types.\nIn some instances, abrasive cleaning surfaces have been attached to certain gloves or mittens, including gloves or mittens which are waterproof. To the knowledge of the inventor, these abrasive or cleaning surfaces have been attached to the gloves adjacent to or at the palm area of the particular glove. The prior art in this respect is explified by U.S. Pat. Nos. 2,459,521 to Woodbury; 3,643,386 to Grizyll; and 4,038,787 to Bianchi; as well as French Pat. No. 2,278,277 to Scott.\nThe above and other art known to the Inventor relates only to the attachment of the abrasive or cleaning material to the palm area. The problem encountered therein is the palm of the user is soft and, therefore, is not as effective for the cleaning of difficult surfaces as is a scouring stone or other material having a rigid abrading surface. Further, in that the area of the palm is relatively large, the pounds per square inch that may be applied therefrom, to the working surface, is relatively small.\nA shortcoming in prior art cleaning gloves has been the lack of physical flexibility, their incapacity to act as a sponge in order to accomplish the absorbtion and retention of detergents in solution, and the single cleaning purpose thereof, i.e., particularity of the cleaning surface for which a given cleaning abrasive glove may be used.\nThe inventor has discovered that the location of cleaning or scouring units upon the outer portion of the glove and, more particularly, upon the second row of phalanges of the middle three fingers of the hand of the user, will make possible the application of many more pounds per square inch of force than is possible in the usage of the palm of the hand as a cleaning or scouring means.\nThe inventor has also discovered a usage, and interrelationship of usage, of certain materials for (a) the glove proper and (b) the cleaning or abrasive material that, in combination, provide enhanced effectiveness in terms of cleaning capability, variety of material that may be cleaned, and comfort to the user.\nThe present invention is believed to be properly classified in one or more of the following areas: U.S. Class 2, Subclass 161; Class 51, Subclass 391; and Class 242, Subclass 15."}
{"text": "The invention relates generally to agricultural vehicles, and, more specifically, to a status indicator for an autonomous agricultural vehicle.\nVehicles, such as agricultural tractors, may be driven through a field of crops to perform various agricultural operations. In recent years, agricultural vehicles have been designed to operate at least partially without input from an onboard operator. For example, the vehicle may perform one or more operations by receiving one or more instructions remotely, thereby enabling the vehicle to operate at least partially in an autonomous manner (e.g., without input from an operator within the vehicle). For example, a remote operator or instructions from a remote control system may activate the vehicle and initiate operations."}
{"text": "a. Field of the Invention\nThe instant invention relates to catheters having an outer surface with multiple aspiration passageways, through which fluids, such as blood, may flow into an interior lumen of the catheter, such as a drainage lumen, an aspiration lumen, or even the central lumen of the catheter. The pattern, design, placement and number of aspiration passageways, or holes, allows for improved catheter performance during procedures, such as transseptal procedures, where the distal tip commonly becomes occluded, while still maintaining the structural integrity and functionality of the catheter.\nb. Background Art\nCatheters are used in a variety of medical procedures to gain access to interior regions of the body and to perform a variety of procedures, such as the ablation of targeted tissue areas. During such procedures, a physician must carefully and precisely control the position of the catheter. To achieve more precise control of instruments such as ablation catheters, various guiding introducer catheters have been developed.\nTypically, in procedures where introducer catheters are used to guide an ablation catheter to a particular location in the heart, a small incision is first made at an appropriate location, for example, the groin, the arm, or any other location known to those skilled in the art, and then an artery or vein is punctured at a relatively shallow angle with an appropriate needle, for example, a stylet or some other puncturing device known to those skilled in the art. The needle is then partially withdrawn from the artery or vein, and reinserted at a steeper angle A guidewire is inserted into the artery or vein and advanced to a desired region, for example to the heart region. An introducer catheter is directed over the guidewire to the desired portion of the heart. An instrument, often embodied in a separate catheter, to be used for the treatment of the heart or the heart region, for example, an ablation catheter, may then be advanced through the lumen of the first catheter and placed at the desired location to perform the desired procedure.\nThe specific type of introducer catheter used will depend upon what portions of the heart the operation is being performed on, and accordingly whether a transseptal or retrograde procedure is used. In transseptal procedures, for example, a catheter apparatus is introduced into the right femoral vein and advanced through the inferior vena cava into the right atrium. A puncture is then made through the fossa ovalis in the inter-atrial septum and the apparatus is advanced into the left atrium.\nOne skilled in the art recognizes the many problems inherent in the use of introducer catheters in such procedures. For example, given the intricate anatomy of the heart, it is possible for the distal end of the introducer catheter to become occluded, particularly if the distal tip becomes pressed against tissue, or in extreme cases, blocked with debris, such as plaque or cholesterol. This problem is particularly pervasive when puncturing the fossa ovalis in the aforementioned transseptal procedure. Furthermore, puncturing the fossa ovalis often causes localized bleeding, which can lead to the coagulation of blood and subsequent clogging of any openings at or near the distal end of the introducer catheter.\nBecause an ablation catheter, when advancing through an interior lumen of the transseptal introducer catheter, is cooled in part by the flow of blood through introducer's interior lumen, an occluded distal end would impede the flow of blood and interfere with the temperature control of the ablation device. Accordingly, it is desirable to place aspiration passageways in the side of the catheters to ensure that blood or other fluids continuously flow through an interior lumen. These aspiration passageways function as relief holes in the side of the introducer catheter which allow for the flow of blood into an interior lumen and through the catheter when the distal end of the catheter is occluded. Therefore, it is one object of the present invention to provide improved catheters for a variety of transseptal and retrograde procedures, where the design of aspiration passageways prevents the occlusion of the catheter."}
{"text": "1. Field of the Invention\nThe invention pertains to the field of hearing protection. More particularly, the invention pertains to hearing protectors having talk-through capability.\n2. Description of Related Art\nTalk-through listening systems typically comprise hearing protection devices (i.e. earplugs, headsets, helmets) with microphones mounted on their outside surface. The microphones convert sound pressure levels (SPLs) to an electrical signal. The electrical signals are amplified and used to drive audio transducers such as speakers in headphones, helmets, or earplugs or the like. Electronic automatic gain control and limiting devices ensure that the SPL generated in the canals of the user are kept within safe levels. In this way, the user can monitor sound in a high-SPL environment at safe SPLs.\nOne example of such an environment is the flight deck of aircraft carriers. Aircraft maintainers need to service aircraft under extremely high noise conditions. However, they need to be aware of their surroundings and must communicate to each other. The maintainers must know if a vehicle or plane is approaching from areas outside of their view. For this, they use acoustic cues. Moreover, the maintainers often need to communicate with each other face-to-face in high acoustic noise, without the use of radios."}
{"text": "This invention relates generally to dual-data rate (DDR) memories, and more particularly to full page increment/decrement burst for such memories.\nGraphics-intensive applications for computers such as personal computers (PC\"\"s) are becoming increasingly more popular. Such applications include high-end computer-aided drafting (CAD) applications, games, real-time video applications, as well as other applications. As these applications become more complex, they require the computers on which they are run to render and execute graphics much more quickly. Furthermore, as the typical resolution of computer screens have increased from 640xc3x97480 pixels (horizontalxc3x97vertical) to 800xc3x97600, 1024xc3x97768 and beyond, and increased color information per pixel from two bits to 24 bits to 32 bits and beyond, the processing demand placed on the computers for fast graphics execution has also grown.\nThe typical computer relies on a graphics card (also known as a video card, graphic accelerator card, or a display adapter, among other terms) to assist it in the display of graphics on a display device. A graphics card generally includes a specialized processor or processors that are tailor-made for graphics rendering, as well as an amount of memory, ranging from one, two, four, eight, sixteen megabytes and up, so that a complete screen of graphics information, known as a frame, can be stored by the graphics card. Thus, this memory is generally known as a frame buffer of the graphics card. Graphics xe2x80x9ccardsxe2x80x9d may also be integrated within a single chip on a motherboard of a computer. A graphics card, and potentially other components, make up the graphics subsystem of a computer.\nInitially, the memory of a graphics card was standard-issue dynamic random-access memory (DRAM), of a sort also used by computer processors to hold more general information. Thus, as improvements in memory to increase their speed became available, such as the introduction of synchronous dynamic random-access memory (SDRAM), they usually have been utilized within graphics cards, too. Ultimately, however, the specialized needs of graphics rendering required their own type of memory, such as synchronous graphics random-access memory (SGRAM), which is analogous to SDRAM, but includes enhanced graphics features for use with graphics cards. The need for faster memory within graphics cards has not, however, abated.\nThus, graphics cards manufacturers have looked to new technologies, such as Rambus DRAM\"\"s (also known as Direct RDRAM\"\"s), available from Rambus, Inc. of Mountain View, Calif., to increase graphics subsystem performance. Rambus DRAM use within graphics cards, however, has been limited because it is based on a closed standard governed by Rambus, Inc., such that use of Rambus DRAM requires the payment of royalties to Rambus, Inc. Therefore, manufacturers have looked to other technologies that are based on open standards.\nOne such type of memory is the Double Data Rate (DDR) DRAM. The DDR DRAM achieves increased performance by providing for two data accesses within a single clock cyclexe2x80x94hence its namexe2x80x94by enabling the memory to read data on both the rising and falling edges of each clock cycle. The concept of DDR memories has been extended to SDRAM\"\"s and SGRAM\"\"s in particular, resulting in DDR SDRAM and DDR SGRAM. Such memory has witnessed increased interest on the part of graphics card designers as a manner by which increased graphics performance can be realized.\nA disadvantage to DDR SDRAM/SGRAM\"\"s found in the prior art, however, is that full-page burst is not generally provided for. (In general, a burst operation for a memory is defined as an operation retrieving a given number of data stored at sequential locations within the memory (e.g., a full-page of memory), which ultimately allows the data to be retrieved in a faster mannerxe2x80x94hence the term xe2x80x9cburst.xe2x80x9d) This is because of a limitation of the prefetch nature of DDR devices. That is, in a DDR device, for a given clock cycle, two words of data, each of n-bits length, are retrieved, such that both must be from the same location within the memory as addressed by the logical circuitry before the next location can be moved to. This is acceptable for a full-page increment burst starting with an even start word address, because the second data word retrieved is still within the same location within the memory as the first data word. However, a full-page increment burst starting with an odd start word address does not work, because the second data word retrieved will necessarily not lie within the same location as the first data word as addressed by the logical circuitry, violating the limitation of the prefetch nature of DDR devices.\nThere is a need, therefore, for a DDR SDRAM/SGRAM that has burst capability. For these and other reasons, there is a need for the present invention.\nThe above-mentioned shortcomings, disadvantages and problems are addressed by the present invention, which will be understood by reading and studying the following specification. The invention relates to a dual-data rate (DDR) synchronous dynamic random access memory (SDRAM)/synchronous graphic random access memory (SGRAM). In one embodiment, the DDR SDRAM/SGRAM includes a memory array and a logic circuitry coupled thereto. The memory array is addressable by even and odd word addresses. The logic circuitry has a burst increment mode to access the array starting at an even word address and a burst decrement mode to access the array starting at an odd word address.\nThus, the invention provides for advantages not found in the prior art. A DDR SDRAM/SGRAM according to the invention is able to have a burst increment mode when starting at an even word address, and a burst decrement mode when starting at an odd word address. Thus, the limitation imposed by the prefetch nature of DDR devices is not violated. When starting at an even word address, the DDR SDRAM/SGRAM counts up (increments) so that the second data word retrieved is still from the same memory location as the first data word as addressed by the logic circuitry, while when starting at an odd word address, the DDR SDRAM/SGRAM counts down (decrements), also so that the second data word retrieved is still from the same memory location as the first data word as addressed by the logic circuitry.\nThe invention includes DDR SDRAM/SGRAM\"\"s, memory devices, graphics cards, computers and methods of varying scope. Still other embodiments, advantages and aspects of the invention will be understood by reference to the drawings and by reading the detailed description."}
{"text": "Insulating glass assemblies for use in the manufacture of windows, doors and the like commonly have two substantially parallel, spaced-apart glass panes spaced apart by a peripheral spacer. Spacers commonly are of metal, usually of tubular configuration, that are formed so as to have two flat, substantially parallel sides facing the confronting surfaces of the panes and bent so as to conform to the periphery of the glass panes. Sealant materials such as polyisobutylene are employed between the flat sides of the spacer and the confronting glass surfaces to seal the glass surfaces to the spacer. To enhance the thermal resistance across the glass assemblies, the interpane space may be filled with an insulating gas such as argon having a thermal conductivity that is less than that of air.\nIn the manufacture of insulating glass units, uniform production line procedures enable glass assemblies of a single size to be made in large quantities. Custom insulating glass units, on the other hand, are generally manufactured in quantities as small as a single unit, and a single order may require the manufacture of units having varying sizes and shapes.\nVarious methods and apparatuses have been suggested to enable air within the interpane space to be replaced with an insulating gas such as argon. In one method, the glass panes are adhered to a spacer to form a substantially sealed interpane space, and then air within the space is gradually replaced with argon through an access port. In another method, the interpane space of a multipane glass assembly is filled with an insulating gas by first drawing a vacuum to remove air from the interpane space before both panes are sealed to the spacer, and then charging the evacuated interpane space with an insulating gas. After the interpane space is filled with the insulating gas, the panes are sealed to the spacer.\nVarious methods and apparatuses for replacing air with an insulating gas in insulating glass units are shown in U.S. Pat. Nos. 5,017,252, 4,780,164, 5,573,618 (Rueckheim) and U.S. Pat. No. 5,476,124 (Lisec). In the last mentioned patent, an apparatus is described in which an insulating glass unit having a pair of glass panes separated by a peripheral spacer is conveyed by a conveyor belt between parallel plates, the bottom edge of the outer glass pane being spaced slightly away from the spacer to provide generally vertical openings along the side edges of the unit. The leading edges of the glass panes are conveyed into contact with a vertical sealing device. Another vertical sealing device is then moved into contact with the trailing edge of the glass panes to seal, with the gas-tight conveyor belt, the space between the glass panes. An insulating gas is then flowed laterally from one vertical sealing device to the other under conditions avoiding turbulence. When the glass unit has been appropriately filled with insulating gas, one plate is advanced toward the other to compress the glass unit between the plates and thus completely adhere the glass panes to the peripheral spacer. This device replaces air with an insulating gas in one glass unit at a time, and due to its employment of non-turbulent gas flow, requires considerable time to replace the air with insulating gas. It would be advantageous to provide a method and apparatus for filling one or a plurality of the same or different size insulating glass units at a time with an insulating gas in a manner providing rapid and substantially complete replacement of air.\nApplicant's U.S. Pat. No. 5,957,169 teaches an apparatus and method of filling insulating glass with insulating gas. In accordance with this method, a partially assembled glass unit is conveyed within an enclosure. This partially assembled glass unit has a pair of spaced panes of glass and a peripheral spacer, with the lower edge of one pane being spaced from the spacer to provide a bottom gap. An insulating gas is introduced under turbulent flow conditions upwardly through the gap in the partially assembled glass unit to turbulently mix with the air present therein. A mixture of insulating gas and air is exhausted from the enclosure until the concentration of insulating gas within the enclosure reaches a desired level. The lower edge of the glass pane of the partially assembled glass unit is then closed against the spacer to seal the interpane space.\nOne problem encountered in manufacturing insulating glass assemblies is that the seal between one or both of the glass panes and the spacer fails. While this can happen over time, this defect more commonly occurs at the time of manufacture due to an improper application of the sealant about the periphery of the frame. In most cases, the sealant is not applied uniformly and relatively small, pin-hole-sized pathways through the seal remain even after the glass panes are urged into contact with the spacer.\nSeal failures can cause a variety of problems in insulating glass units. For example, a seal imperfection will cause a leak of the insulating gas from the interpane space, and the pressure in the glass unit may decrease to a point where the pressure is no longer sufficient to support the upper pane, and this in turn may cause the upper pane of glass of the unit to bow or sag under its own weight. Moreover, such a leak can allow the introduction of moisture and other unwanted environmental elements, leading to a faulty glass unit. A leak in the seal can thus lead to eventual failure of the entire glass unit, resulting in the most extreme cases in a blowout of the glass panes and/or the spacer and seal.\nIt has proven difficult to reliably detect imperfections or failures in the peripheral seal provided by the spacer. Applicant has found that the upper pane of glass in larger units with inadequate seals will tend to sag slightly if the unit is laid with the panes oriented horizontally. It has been found useful in commercial production for over a year to fill the insulating glass unit with an insulating gas at roughly atmospheric pressure and profile the shape of the upper sheet of glass. The upper sheet of glass will tend to sag slightly in the center under its own weight, but the insulating gas in the interpane space of a properly sealed unit will help support the glass and limit its deflection. If the seal has a significant leak, though, the outer surface of the glass unit will tend to be concave. Such a concave profile indicates that insulating gas is being forced out of the interpane space under the weight of the upper glass pane, thereby lowering the pressure of gas within the glass unit and allowing the upper glass pane to sag under its own weight.\nUnfortunately, this process is not a reliable indicator of seal failure in smaller glass units. In larger glass units, the center of the glass is spaced sufficiently far away from the spacer to permit measurable deflection. In smaller units, however, the distance between the supported edge of the upper pane and its center of mass is smaller. As a consequence, the bending moment on the glass often falls short of that necessary to consistently induce a sagging of the glass pane significant enough to fall outside standard manufacturing tolerances.\nHence, monitoring thickness of larger panes can serve as a reasonably reliable means for detecting inadequate seals in insulating glass units. In smaller units, though, this measure is not a strong, statistically reliable indicator of seal failure. It would be advantageous to provide a method for reliably detecting seal imperfections or failures in insulating glass units of varying sizes and shapes, in a manner providing rapid detection of imperfection or failure regardless of the size of the glass unit."}
{"text": "1. Field of the Invention\nThe present invention relates to semiconductor memory devices. More particularly, the present invention relates to a repair circuit and a method of repairing defects in a semiconductor memory device.\nThis application claims priority under 35 USC § 119 to Korean Patent Application No. 2005-117258, filed on Dec. 3, 2005, the subject matter of which is incorporated herein in its entirety by reference.\n2. Description of the Related Art\nA semiconductor memory device generally includes a main memory array having a plurality of memory cells, and a redundant memory array having redundant memory cells adapted for use in the repair of defective memory cells in the main array. It is preferable that all memory cells of the main memory array operate normally. However, defective memory cells are often formed within the main memory array due to variations in the manufacturing processes used to fabricate the memory cells as well as variations in electrical performance characteristics across a large number of memory cells. When a defective memory cell is identified in the main memory array, it may be replaced (or substituted) with one of the redundant memory cells. This capability ensures uniformly acceptable memory cell performance across the entire main memory array.\nIn general, at various stages during the lengthy sequence of fabrication processes used to manufacture a semiconductor memory device, tests are performed to determine whether or not the memory cells are operating normally. After each test, a memory cell identified as defective may be replaced with a redundant memory cell. This test, identification, and replacement procedure is typically performed while the semiconductor memory device is still in its “wafer state”, i.e., before individual dies are cut from the fabrication wafer and packaged into individual device form.\nThere are many different types of memory cell defects, including row related defects and/or column related defects. That is, while some memory cell defects are identified in relation to a particular memory cell, memory cell defects are more commonly identified as a defective row or column. As a result, semiconductor memory devices require both a row repair circuit and a column repair circuit in order to remedy row defects and column defects, respectively.\nConventional repair techniques may be divided into before-package repair (BPR) techniques—adapted to repair identified defects in wafer state devices, and post-package repair (PPR) techniques—adapted to repair defects in a packaged devices. PPR techniques are commonly used to repair defective memory cells which remained unidentified in the wafer state, but which were subsequently identified during package level testing. Both PPR and BPR techniques require row repair and column repair circuits.\nSome PPR techniques use electrical fuse circuits (e-fuses) instead of laser fuse circuits to repair defective memory cells and improve package yield for the semiconductor memory devices. E-fuse circuits offer several performance advantages, but tend to be more complicated in their layout relative to laser fuse circuits. Thus, conventional PPR techniques incorporating both the row and column repair circuits using e-fuse circuits occupy an unacceptably large layout area within the overall semiconductor memory device. This quality is at odds with continuing pressure to increase the fabrication density of memory devices."}
{"text": "This invention relates to a process for the preparation of polyurethanes from diisocyanates containing uretdione rings.\nIt is known in the art (DE-PS No. 1,153,900) to react diisocyanates containing uretdione groups with relatively high molecular weight polyhydroxyl compounds and optionally difunctional chain lengthening agents to produce polyurethanes. In order to achieve sufficient polyaddition of the uretdione-containing diisocyanates (which are generally high melting and only sparingly soluble), it is usually necessary to employ reaction temperatures above 100.degree. C. Use of such temperatures may however result in the opening of the uretdione rings and the initially straight chained polyurethanes may undergo a much greater degree of cross-linking than desired.\nSubstantial lowering of the solidification temperature and shortening of the solidification time in the curing process by using catalysts of the kind conventionally used in polyurethane chemistry, e.g. organic metal salts or tertiary amines have been attempted. However, the density of cross-linking is found to be undesirably increased and the mechanical properties of the polyurethane are seriously impaired.\nThe same problems are encountered when solid uretdione diisocyanates with delayed reactivity such as those which have been described in DE-OS No. 3,230,757 are used for the preparation of storage stable mixtures of so-called one component systems for the production of polyurethanes. In these cases, again, the use of organic lead or tin compounds as the preferred catalysts cannot completely prevent opening of the uretdione rings and the attendant undesirable side reactions."}
{"text": "Successful marketing of products (including goods and/or services) relies upon many factors; however, a predominant factor is word of mouth, or the effect of interpersonal communication. Word of mouth marketing relies upon the recommendation or opinion of a trusted colleague, friends, or relatives of the consumer. Word of mouth marketing is effective yet is difficult for merchants and brand owners to successfully perform.\nVarious attempts have been made to determine consumer social networks, and hence to promote or influence word of mouth marketing through particularly influential members of such networks. However, these attempts have generally not been successful, despite (or rather because of) the massive amount of available data. Thus, merchants frequently attempt instead to analyze and predict the behavior of each consumer individually, and hence disregard such connections.\nThere are many efforts organizations employ to use the power of word of mouth.\nUndercover Marketing—\ncommercially motivate people to recommend a certain product (personally or through internet chat rooms, talkbacks etc.) or use the product or service in public.\nCreating Buzz—\nCertain companies use word of mouth advertising agencies to identify opinion leaders and use them to spread word of mouth. These opinion leaders are recruited one by one through web site or personal interviews.\nTargeting Trend Setters—\nCertain companies send a product to several A-list trend setters to spread the word. For example, the book “the da Vinci code” was sent to 10,000 industry trend setters for them to create and amplify the word of mouth around the book.\nThus, none of the above attempted solutions successfully uses social networks to determine how to target word of mouth marketing to influential members of such networks."}
{"text": "Mounting and demounting tires onto and from a wheel can be a daunting task. The beads of a tire are designed to fit snuggly to the inner side of the flange of a wheel. Stretching the bead a sufficient amount to clear the flange is difficult, particularly when working with larger wheels and tires such as in the case of those for a truck. The task is made all the more difficult when the wheel is made of a soft material that can easily scratch or crack, such as aluminum. Thus, in this case extra care in an already difficult job is mandatory.\nThere have been many advances proposed in the field to make the job of working a tire easier; these working tools are well known in the art. Generally, these devices consist of a lever arm for gaining leverage, the lever arm having at least one working end configured to receive and work a bead of the tire. Early examples of such devices are seen in U.S. Pat. No. 834,908 to Hussey and U.S. Pat. No. 951,200 to Pilliner, each of which teach a tool for detaching and resetting tires. A more recent example of a typical device is seen in U.S. Pat. No. 6,508,479 to Kliskey. This device teaches a tool for demounting a tire, the tool consisting of a handle and working ends protruding generally perpendicular from the handle, the working ends adapted to work a first tire bead and a second tire bead, respectively.\nFurther examples of inventive devices is in the field include U.S. Pat. Nos. 1,100,032, 1,567,025, 1,741,801, 1,829,804, 2,112,661, 2,188,211, 2,226,757, 2,294,271, 2,311,789, 2,344,704, 2,399,146, 2,571,517, 2,643,803, 2,712,350, 3,164,197, 3,584,672, 3,908,728, 4,436,134, 4,527,607, 4,919,184, 5,123,470, 5,143,134, 5,213,146, 5,265,661, 5,343,921, 5,363,897, 5,417,270, 5,740,848, 5,806,578, 5,971,052, 6,179,032, 6,182,735, 6,192,959, 6,273,171, 6,712,114, D134,799, D142,931, D158,110, D158,629, D264,043, D335,754, D395,809, D412,271, and D476,205.\nA proposed advantage of many of the aforementioned devices is that in addition to offering leverage to mount or demount a tire, they provide a rim or flange catching feature to prevent the device or tire from slipping out of position when the tire is being worked.\nA shortcoming of the prior art devices is it that due to the narrow nature of the working end features, the tools do not offer optimum leverage since they do not spread the torque energy of the lever arm beyond the narrow working end. A further shortcoming of the prior art is that they tend to twist about the narrow working ends. This can cause the tool to slip out of position when the tire is being worked, thereby, at best, causing a need for the user to start the process over or, at worst, causing injury to the user. This is particular troublesome when attempting to work the first tire bead.\nPast attempts by the inventor to address these shortcomings included the development of a tire working tool consisting of a lever arm with a first working end adapted to accept a clevis, the clevis designed to work the first tire bead over the first flange of a wheel rim. The inventor soon discovered, however, that the design of the clevis was not well suited to the task at hand. The original clevis design was too narrow to provide optimal assistance in working a tire bead. The wings on the clevis neither extended a sufficient distance to prevent the clevis from slipping out of position nor prevented tires from sticking to the wheel rim. Moreover, when working a tire bead, the clevis did not maximize the leveraging energy along the bead of the tire and did not, therefore, reduce the effort required of the user to the degree desired.\nAnother shortcoming of the prior art, including that of the inventor, is that they do not teach or contemplate the added difficulties in working with soft material wheel rims, such as aluminum.\nThere is need, therefore, for a tire working tool that minimizes the risk of injury to a user and damage to tire rims and that, further, makes the demounting and remounting of a tire onto a wheel easier to accomplish than heretofore known. The present invention provides such a tool."}
{"text": "To deal with challenges of wireless broadband technologies, and maintain the leading edge of the 3rd Generation Partnership Project (3GPP) network, the 3GPP has formulated a Long Term Evolution (LTE) plan of a mobile communications network. Under the guidance of the LTE plan, the 3GPP has defined a new mobile communications network architecture, that is, an evolved packet system (EPS) wireless network architecture.\nFIG. 1 is a schematic diagram of an EPS wireless network architecture in the prior art. As shown in FIG. 1, when user equipment (UE) accesses an EPS network by using a serving GPRS support node (SGSN), the UE may access a serving gateway (SGW) by using an S4 interface, and then, access a packet data network gateway (PDN-GW or PGW) by using an S5 interface. Functions of main network entities in the EPS network are described as follows:\nAn evolved universal terrestrial radio access network (E-UTRAN) is a network including multiple evolved NodeBs (eNodeB), and implements a wireless physical layer function, resource scheduling and radio resource management, and a radio access control and mobility management function. The eNodeBs are connected to an SGW by using a user plane interface (S1-U interface), and are configured to transmit user data; and are connected to a mobility management entity (MME) by using a control plane interface (S1-MME interface), and are configured to implement a radio access bearer control function by using an S1 Application Protocol (S1-AP).\nThe MME is mainly responsible for functions such as mobility management of the UE, session management, encryption and integrity protection of non-access stratum (NAS) signaling, management of a tracking area list (Tracking Area List), and PGW/SGW selection.\nAn SGW is mainly responsible for data transmission, forwarding, and route switching of the UE, and is used as a local mobility anchor when the UE is handed over between eNodeBs.\nA PGW is an entry for an external network to send data to the UE, and is mainly responsible for IP address assignment of the UE, data packet filtering, rate control, charging information generation, and the like.\nA policy and charging rules function (PCRF) entity determines a corresponding policy according to a network access limitation of a user, an operator policy, user subscription data, and information about a service that is being performed by the user, and provides the policy to a transmission gateway for execution, to implement policy and charging control.\nOperator's IP services (Operator's IP Services) are implemented in the LTE network by using an IP multimedia subsystem (IMS) network. In addition, a packet switched streaming service (PSS) technology is a technology defined by the 3GPP and used to provide a streaming media service to a user. A PSS network architecture mainly includes a mobile terminal and a PSS server on a network side.\nData packets sent and received by the UE are transmitted in the EPS network by using an EPS bearer. For the convenience of description, the EPS bearer may also be referred to as a bearer in the following specification. Each UE may have multiple bearers, and different bearers can satisfy quality of service (QoS) requirements of different services. An eNodeB and an SGW may store information about each bearer, that is, a bearer context. The bearer context includes an SGW tunneling endpoint identifier (TEID) and an eNodeB TEID of each bearer. The SGW TEID is used for an uplink packet sent by the eNodeB to the SGW and the eNodeB TEID is used for a downlink packet sent by the SGW to the eNodeB. The eNodeB implements bearer context synchronization with an MME by using an S1-AP message, and the SGW implements bearer context synchronization with the MME by using a GPRS tunneling protocol-control plane (GTP-C) message, thereby implementing bearer context synchronization between the eNodeB and the SGW.\nWhen receiving an uplink packet sent by the UE, the eNodeB may encapsulate the uplink packet of the UE into a GTP-U packet according to the bearer context. The GTP-U packet includes a GTP-U header, and the GTP-U header includes an SGW TEID. Because different bearers may use different SGW TEIDs, when receiving the GTP-U packet sent by the eNodeB, the SGW may determine, according to the GTP-U header, a bearer to which the GTP-U packet belongs. Similarly, when receiving a downlink packet sent to the UE, the SGW encapsulates the packet of the UE into a GTP-U packet. The GTP-U packet includes a GTP-U header, and the GTP-U header includes an eNodeB TEID. Therefore, when receiving the GTP-U packet sent by the SGW, the eNodeB can determine, according to the GTP-U header, a bearer to which the GTP-U packet belongs.\nA mobile edge computing (MEC) platform is a logic network element deployed on the S1-U interface, and is located between the eNodeB and the SGW. The MEC platform mainly includes a traffic bus (traffic bus) and an application (application). The traffic bus is responsible for obtaining a data packet of the UE and forwarding the data packet to a corresponding application. After processing the data packet, the application sends the processed data packet to the traffic bus for routing.\nMultiple types of applications may be installed on the MEC platform, to improve service experience of the user. The application on the MEC platform may capture an uplink packet sent by the UE, and modify, check, and forward the uplink packet of the UE, or may directly reply to the uplink packet sent by the UE. For example, a video buffering application may be installed on the MEC platform. When the UE sends a video request packet, the video buffering application processes the video request packet of the UE. If the video buffering application does not include a video requested by the UE, the video buffering application may forward the video request packet of the user to the SGW, or if the video buffering application stores a video requested by the UE, the video buffering application directly sends a video reply packet to the UE.\nThe applicants find that when the MEC platform directly replies to the uplink packet of the UE, the MEC platform cannot correctly perform GTP-U encapsulation, and consequently, cannot send a reply packet corresponding to the uplink packet of the UE to the eNodeB."}
{"text": "Conventionally, when print jobs are executed continuously in a printer, calibration is executed in the middle of a job or between jobs, interruption of printing due to the execution of the calibration is controlled as much as possible, and the calibration is executed after the jobs are ended, to appropriately prevent the degradation of printing characteristics."}
{"text": "The present invention relates to improvement of an apparatus for making baked goods such as cakes, bread and bun.\nThe conventional baking equipment used in a bakery or confectionary factory to bake bread, bun or cakes includes a conveyor 20 to 60 meters long on which cake and bread pans are transported. The bread, cakes and bun are baked by heat sources arranged along the conveyor. This eqipment is very large and occupies a large floor space."}
{"text": "It has long been a primary objective of those working in the filing cabinet and business office furniture art to provide means whereby a filing cabinet may not only be locked against unauthorized entry but also to provide an anti-tip mechanism which secures the doors and/or drawers of the cabinet against the sequential opening of more than one drawer at any one particular time. It has also been a further and long time objective of those versed in the art to find a means whereby the foregoing objective could be accomplished in economical fashion. That these objectives have not been satisfactorily achieved is evidenced by the multitude of inventions proposed to this end, representative of which are those disclosed in the following U.S. Pat. Nos. 3,900,236 issued Aug. 19, 1975; 3,199,937 issued Aug. 10, 1965; 3,888,558 issued June 10, 1975; 3,404,929 issued Oct. 8, 1968; 3,883,220 issued May 13, 1975; 3,602,564 issued Aug. 31, 1971; and 3,297,376 issued Jan. 10, 1967. These patents may be relevant to the invention claimed herein in that they relate to filing cabinets having anti-tip structure.\nFiling cabinets heretofore proposed, such as those patents listed above, have suffered from a variety of deficiencies and difficulties. While the prior art anti-tip mechanisms have been generally successful in preventing more than one drawer from being opened at a time when the drawers are opened only one at a time, most such anti-tip mechanisms have been ineffective in preventing more than one drawer from being opened simultaneously. Thus, U.S. Pat. Nos. 3,888,558, 3,883,200 and 3,404,929 are examples of cabinets which do not prevent the simultaneous opening of more than one drawer and U.S. Pat. No. 3,602,564 is exemplary of a cabinet which includes a door interlock which does not prevent more than one door from being opened simultaneously. This deficiency in the prior art designs may lead to the upsetting of the entire cabinet when the cabinet itself is tipped slightly so that the drawers tend to run open of their own weight at the same time: a not uncommon circumstance when the cabinet is being moved. The solution to this problem proposed by U.S. Pat. Nos. 3,900,236 and 3,199,937 involves significant complication of the design of the anti-tip mechanism, which complications lead to excessive costs and difficulties in assembly of the cabinets.\nAn additional problem which has been encountered when prior anti-tip devices have been incorporated in filing cabinets is that such mechanisms frequently add additional depth to the cabinet and/or do not permit the drawer of the cabinet to be opened to a degree that the contents of the drawer can be removed vertically from the opened drawer. This difficulty is magnified when the cabinet is of the lateral file type in which cabinet depth is of primary importance and in which vertical deposit and removal of the filing material is a requirement. In an attempt to address this requirement, U.S. Pat. No. 3,900,236 has proposed a device having a complicated camming keeper mechanism which is subject to failure in the event that debris or rust were to interfere with the motion of the cam wheel and is both expensive and difficult to fabricate. Furthermore, since the keeper mechanism is actuated upon opening and removal of one drawer from the cabinet, and is not deactivated until the drawer is returned to its closed position, all of the remaining drawers are prevented from complete removal from the cabinet in the event that it is desirable to lighten the cabinet to facilitate its transport by the complete removal of all of the drawers. Accordingly, the file cabinet disclosed by U.S. Pat. No. 3,900,236 can be lightened by the complete removal of only one drawer at any one time, which capability may not be of any significant help when the cabinet is transported. Further lightening of the cabinet requires the inconvenent expedient of removing the contents of the other drawers rather than complete removal of the drawers themselves.\nThe foregoing, as well as other problems encountered by prior art devices have been generally overcome by my prior U.S. Pat. No. 4,272,138 issued on June 9, 1981. In this patent the anti-tip mechanism includes a segmented column consisting of a plurality of longitudinally extending locking bars or snubber elements axially aligned in ent-to-end abutting relation. The columns extend adjacent to each one of a plurality of horizontally movable components such as roll-out shelves, drawers and/or other suitable storage units together with an equal number of doors movable from a closed, blocking position to an open, releasing position. The columns are each resiliently mounted within the cabinet by a plurality of leaf-type springs which function to bias the columns to a central equilibrium position and to permit movement thereof to one of two axially displaced positions. Each drawer carries on its side an activator rail which is disposed at right angles to the column. The tapered end of each activator rail is positioned so as to insinuate itself between two adjacent abutting snubber elements when the respective component is opened. The act of activator rail insulation between adjacent snubber elements causes each of the snubber elements in the column to be displaced from a central equilibrium position to one of two displaced positions. In their displaced positions, the snubber elements block subsequent opening of any second component.\nAdditionally, the cabinet of my aforesaid patent includes a lock means for simultaneously locking the doors and the components against movement from their closed to their open positions when in a locking mode and for permitting movement of a single drawer from its closed to its open position while simultaneously locking both the remaining unopened doors and all of the other drawers against movement from their closed to their open position when in an unlocking mode. Such lock means includes a reciprocating key operated plunger lock, a linkage in the form of an angle iron and a pin laterally extending from the uppermost column. The angle iron is disposed horizontally within the cabinet, is pivotally mounted on the end of the plunger lock and is positioned adjacent the laterally extending pin whereby depression of the plunger lock causes the angle iron to pivot into engagement with the pin to cause the segmented locking bars to be driven downwardly into a locking position."}
{"text": "When combined with a lens formed of an ultra-low dispersion glass, a lens formed of a high-refractivity low-dispersion glass enables the downsizing of an optical system while correcting chromatic aberration. It hence occupies a very important place as an optical element to constitute an image-sensing optical system or a projection optical system such as a projector.\nPatent Document 1 discloses such a high-refractivity low-dispersion glass."}
{"text": "1. Field of the Invention\nThis earth soil sampler is specifically adapted to extract and contain successive samples of earth soil such as may have been contaminated by oil spillage. For example, earth soil in a given ten foot square vicinity of transformers, may be contaminated by oil containing PCB and if this likelihood occurs, it is desirable to secure small samples which cumulatively comprise a 5 gram total. This total, in practice, comprises 10 or more probing extractions and is desirably deposited in a 20 ml glass container for testing.\n2. Description of the Prior Art\nNumerous means have been created to obtain small earth samplings For example, operators have secured a sampler tube in extension of a rod. In this method the operator would hold the rod and while walking over the ground, probe with the tube end numerous times to obtain a cumulative sample. The resultant disadvantageous gravitational effect of packing the tube precludes successive sampling of minute aliquots in a given operation. Furthermore, such samples as were obtained in this manner became difficult to remove. This development had evolved from the classic mode of obtaining samples through the use of a spoon or small trowel wherein plural samples were scraped from the surface of the suspected contaminated soil in a 10 foot square or more area. In this earlier method, the operator desirably estimated his individual samplings to be at approximately 0.5 grams each and in digging over a given area, he would secure upwards of 10-12 samples, adding each to the contents of the same sample bottle. Not only in the art was it difficult for the operator to thus follow a definite pattern such as will insure a relatively complete coverage, but also it was difficult to obtain properly representative examples using such procedures. It was moreover difficult for operators to produce consistently sized sub-samples.\nOther known prior art is present in such devices as may be used for aerating turf. There are also certain pertinent soil testing and coring devices. The best known aerators are represented by Schell et al U.S. Pat. No. 3,163,456 and Jones Pat. No. 522,286 (British) as well as Broadbent Pat. No. 856,537 (British). In none of this art, is the problem addressed of securing minute successive samples for instantaneous containment in a frangible glass bottle, the contents of which are sequentially adapted to testing. The inexact art or aerating thus does not address the more exact problem of soil sampling such as is addressed by the present invention. Coring devices are represented by U.S. Pat. No. 1,109,446 to Melberg and U.S. Pat. No. 2,666,330 to McAndrew. In the Melberg patent a so-called soil magazine is incorporated into the interior of a lightweight shaft for the purpose of obtaining a given soil sample. No effort is made to restrict the in-feed for a given sample nor to effectively cushion the soil magazine against fracture, considering the relative mass of the shaft, as in the present invention. The McAndrew coring device while including a container within, makes no provisions for the controlled obtaining of minute samples in a frangible container such as in the present invention."}
{"text": "In a communication process, trunk line butting interconnection errors (including cross-connected lines and loopback) in TUP (Telephone User Part) or ISUP (Integrated Services Digital Network User Part) results in cross talk, mono-directional connect, and bidirectional disconnect, etc. in sessions. Usually, the trunk errors need to be located with methods such as loopback, pulling E1, or specifying a circuit for Call-Testing.\nHowever, in TDM (Time Division Multiplexing) A-interface circuits, errors may only be checked with methods including loopback, pulling E1 and specifying a circuit for a Call-Testing. Suppose that there are 10 E1 lines between two NEs (network elements), the 10 E1 interfaces (No. 1-10) of NE A need to connect to the 10 E1 interfaces (No. 1-10) of NE B; first, the person responsible for the maintenance of NE B performs “loopback” or “pulling/plugging E1” for E1 interface 1, making the E1 connection “disconnected” or “connected”; then, the person responsible for the maintenance of NE A checks whether a corresponding E1 “disconnected” or “connected” occurs at E1 interface 1 of NE A, to determine whether the E1 connection is correct. All E1 ports are checked by repeating the process. If the “specifying a circuit for Call-Testing” method is used, the person responsible for maintenance may perform actual Call-Testing after E1 is connected, and make the call occupy the channel for testing E1 by the function of “specifying Call-Testing; and if the session is normal, it shows that the E1 connection is correct.\nIt is apparent that the above solutions are troublesome and may not locate circuit errors quickly and conveniently; therefore, the maintenance cost is increased."}
{"text": "A cluster is a set of computers that appear to be a single host computer when viewed from outside the cluster. The computers are separate processors, or nodes, which are interconnected and utilized as a single, unified computing resource.\nA collection of several server computer nodes into a single, unified cluster makes it possible to share a computing load without users or clients needing to know that more than one server is involved. For example, if any resource in the cluster fails, the cluster, as a whole, may continue to offer service to clients using a resource on one of the other servers in the cluster, regardless of whether the failed component is a hardware or software resource.\nVarious software products are available for managing a clustered environment that includes at least two nodes. One such product, offered by Microsoft Corporation, is called Microsoft Cluster Server (MSCS). MSCS improves resource availability by allowing applications and resources to “fail-over” from a server that is not functioning correctly to another server in the cluster.\nIn MSCS, a cluster includes a group of nodes organized so that, if a node fails, another node may take over the failed node's resources. Communication between the nodes of the cluster enables the cluster service to detect node failures. The cluster service keeps track of the current state of the nodes within a cluster (i.e., active node or inactive node) and determines when a node and its resources should fail-over to an alternate node (after a fault of the resource to be failed-over).\nIn MSCS, a node typically includes a computer system running MS Windows™ NT Server, MS Windows™ 2000 Server, or MS Windows™ Net Server that is organized into a cluster. Each node manages its own groups of resources. Each group of resources runs together on the same node. When a node fails, all the resources in a group may be moved over to another node in the cluster. Such a group is also known as a fail-over group, because all the resources of the group, namely all software and hardware entities used for processing, fail-over to another node. Examples of resources are: Software Service, Software Program, Network Name, IP Address, Disk Drive, Data Base, etc.\nOne resource that may be included on each node of a cluster is a database management system (DBMS). A DBMS is typically a set of software programs that controls the organization, storage, and retrieval of data included in a database. The DBMS also controls the security and integrity of the database.\nIt is becoming increasingly desirable that DBMS environments should be available 24 hours a day, 7 days a week. This is true for both on-line and batch operations. As such, on-line transaction processing (OLTP) transactions desirably co-exists with traditional batch transactions, such as queued inserts and range updates/deletes. Therefore, it is desirable that the DBMS be constantly available, regardless of failures or faults related to the DBMS running on the cluster.\nMSCS includes a health check facility known as “Looks-Alive.” “Looks-Alive” is a polling facility that executes a regular probe (heartbeat) over the requisite network. The heartbeats are sent from one node in the cluster to its counterpart nodes to detect node failures. Therefore, “Looks-Alive” checks to see if the physical server or node at an opposing side of a cluster is active (i.e., able to respond).\nIn the event that one node in the MSCS cluster detects a communication failure with another node (via “Looks-Alive”), it broadcasts a message to the entire cluster causing all members to verify their view of the current cluster membership. This is called a regroup event. If the node manager on an individual node does not respond, the node is removed from the cluster and its active resource groups are moved to another active node.\nOn a two-node cluster, the node manager moves resource groups from a failed node to the surviving node. On a three- or four node cluster, the node manager may move the resource groups from the failed node to a next node. For example, the resource groups from failed node 2 may be moved in their entirety to node 3. The resource groups from failed node 3 may be moved in their entirety to node 4, and so on.\nAs such, MSCS includes a mechanism for determining if a DBMS on a node in a cluster has failed. Also, MSCS includes a fail-over mechanism for moving the resource groups from a failed node in the cluster to a surviving node in the cluster. In the context of a DBMS, this means that MSCS will restart the DBMS on a surviving node in the cluster in the event of a failure of the DBMS on another node in the cluster.\nUnfortunately, the health check facility included in MSCS only provides the type of reactive health check described above by reference to “Looks Alive.” In other words, the DBMS on a given node may have already failed by the time MSCS detects the problem. As such, critical information in the DBMS may be lost. Further, DBMS resources may be unavailable for a certain period of time during detection of the failure and the subsequent re-start.\nAs such, it would be desirable to have a health check facility included in a DBMS running in a clustered environment, that reduces the potential loss of information and unnecessary unavailability associated with a node or resource failure."}
{"text": "The present invention relates to techniques for providing interoperability between and among disparate entities, applications and services in a network environment. More specifically, embodiments of the invention provide policy management techniques which facilitate such interoperability despite differing policies associated with the different entities, applications services.\nNetwork communications between different enterprises are typically handled using ad hoc solutions. That is, if two enterprises want to communicate efficiently over the Internet, one will typically have to conform to the requirements of the other. Such requirements may relate, for example, to network communication protocols (e.g., FTP vs. HTTP), or to higher level policy issues (e.g., encryption).\nWS-Policy allows providers of web services to define policies for anyone wanting access to their services. A web service provider publishes a WS-Policy description of a service (e.g., on its web site) so that potential consumers of the service can determine whether they are going to be able to consume the service or not. Thus, consumers of such web services are forced to conform to the policies of the service provider, regardless of their own technology or policy constraints. Given the number of partners or customers with which the typical enterprise would like or is required to communicate, such ad hoc approaches are simply no longer practicable."}
{"text": "Vitamin and mineral supplements for human and veterinary use are commonplace. Some diets, heavy physical exercise and disease conditions may require the intake of considerable quantities of minerals apart from those generally obtained through what otherwise would be considered a normal diet. Calcium and trace mineral supplementation is important primarily for those who have inadequate diets, including growing children. Older adults have an additional need for calcium to help prevent the bone loss which occurs as a normal consequence of the aging process. In particular, postmenopausal women need additional calcium due to hormonal changes which can accelerate the bone loss rate leading to a further diminishment in bone mass.\nThe trace minerals which affect bone growth are copper, zinc and manganese. Supplementation of the diet with these minerals along with a highly bioavailable source of calcium is highly desirable. Commercially available mineral supplements are useful in many circumstances where increased mineral intake is desirable. Most of these multi-vitamin and multi-mineral tablets are low in calcium, requiring separate supplementation with calcium sources. In addition, not all calcium sources are equal in terms of bioavailability and absorption. It would be more convenient if all of the minerals could be administered conjointly in a convenient and/or pleasant tasting form which would not require extra attention, planning and implementation by the user. This could be done in the form of foods and beverages as well as in the form of tablets.\nThere are well-recognized problems associated with adding both mineral and calcium supplements to foods and beverages. Some of these are taste; calcium tends to be chalky in flavor. In addition, the solubility of many calcium sources prevents them from being added to many beverages. Others are interactions of the minerals with the food or beverage which affects the stability and/or the bioavailabilty of the product. This invention provides a means for making such product.\nThis invention also relates to methods of building bone in humans and other animals, i.e., for the treatment of age-related bone loss and related disorders. In particular, this invention relates to such methods of treatment by administration of certain calcium salts and the minerals, copper, zinc and manganese.\nCalcium is the fifth most abundant element in the human body. It plays an important role in many physiological processes, including nerve and muscle functions. Not surprisingly, nutritional and metabolic deficiencies of calcium can have broad-ranging adverse effects. Since about 98% to 99% of the body's calcium is found in bone tissues, many of these adverse effects are manifested through deficiencies in the structure, function and integrity of the skeletal system.\nThe most common metabolic bone disorder is osteoporosis. Osteoporosis can be generally defined as the reduction in the quantity of bone, either from the reduction in bone formation or the acceleration of bone resorption, in either event the result is a decrease in the amount of skeletal tissue. In general, there are two types of osteoporosis: primary and secondary. \"Secondary osteoporosis\" is the result of an identifiable disease process or agent. However, approximately 90% of all osteoporosis cases are idiopathic \"primary osteoporosis\". Such primary osteoporosis includes postmenopausal osteoporosis, age-associated osteoporosis (affecting a majority of individuals over the age of 70 to 80), and idiopathic osteoporosis affecting middle-aged and younger men and women.\nFor some osteoporotic individuals the loss of bone tissue is sufficiently great so as to cause mechanical failure of the bone structure. Bone fractures often occur, for example, in the wrist, hip and spine of women suffering from postmenopausal osteoporosis. Kyphosis (abnormally increased curvature of the thoracic spine) may also result.\nThe mechanism of bone loss in osteoporotics is believed to involve an imbalance in the process of \"bone remodeling\". Bone remodeling occurs throughout life, renewing the skeleton and maintaining the strength of bone. Two reactions are involved, bone loss or resorption and bone growth or accretion. This remodeling occurs in a series of discrete pockets of activity in the bone, called \"osteoclasts\" and \"osteoblasts\". Osteoclasts (bone dissolving or resorbing cells) are responsible for the resorption of a portion of bone within the bone matrix, during the resorption process. After resorption, the osteoclasts are followed by the appearance of osteoblasts (bone forming cells), which then refill the resorbed portion with new bone.\nIn young healthy adults, the rate at which the osteoclasts and osteoblasts are formed maintains a balance of bone resorption and bone formation. However, as normal consequency of aging an imbalance in this remodeling process develops, resulting in loss of bone at a rate faster than the accretion of bone. As imbalance continues over time the reduction in bone mass and thus bone strength leads to fractures.\nMany compositions and methods are described in the medical literature for the \"treatment\" of osteoporosis. See, for example, R. C. Haynes, Jr. et al., \"Agents affecting Calcification\", The Pharmacological Basis of Therapeutics, 7th Edition (A. G. Gilman, L. S. Goodman et al., Editors, 1985); and G. D. Whedon et al., \"An Analysis of Current Concepts and Research Interest in Osteoporosis\", Current Advances in Skeletogenesis (A. Ornoy et al., Editors, 1985). Estrogen is often used to affect the metabolism of calcium by influencing the osteoblast cells. Treatments using fluoride have also been described. However, the utility of such agents may be limited, because of possible adverse side effects. See W. A. Peck, et al., Physician's Resource Manual on Osteoporosis (1987), published by the National Osteoporosis Foundation.\nNutritional therapies for osteoporosis have also been proposed. Many calcium-containing compounds and compositions have been described for use as nutritional supplements. Many commercial preparations are also available, typically containing calcium carbonate or calcium phosphate. Other calcium salts have also been described for use in calcium supplements, including calcium lactate, calcium citrate and calcium gluconate.\nU.S. Pat. No. 3,949,098 issued Bangert (assigned Nabisco, 1976) describes a nutritious orange drink concentrate that contains whey protein. The patent suggests the addition of minor amounts of vitamins and other nutrients which include various cupric salts, manganese salts, zinc salts, as well as calcium salts.\nGerman OLS 2,845,570 issued to E.R.E. (Europe Representation Establishment, 1980) describes a honey containing composition. Honey contains low levels of calcium, manganese, copper as well as trace amounts of magnesium, iron, phosphorous, silicon and nickel. The value of honey as a medicant is undisputed according to this patent application. This application claims a honey containing composition with levarotatory ascorbic acid and citric acid. This patent has issued as U.S. Pat. No. 4,243,794 (1981).\nU.S. Pat. No. 4,497,800 issued to Larsen et al (assigned Mead Johnson & Company, 1985) describes a nutritionally complete ready-to-use liquid diet for providing total patient nourishment. The diet contains free amino acids and small peptides, a carbohydrate source, and nutritionally significant amounts of all essential vitamins and minerals, and stabilizers. The minerals include calcium, copper, zinc and manganese, among others. Most of these minerals are given as the gluconate salt.\n\"Effects of calcium carbonate in hydroxyapatite on zinc and iron retention in postmenopausal women\", Dawson-Hughes, Seligson and Hughes, American Journal of Clinical Nutrition, 44, 83-88 (1986) describes the effect of calcium carbonate on whole-body retention of zinc and iron in thirteen healthy post menopausal women. The test meal, including both dry food and a formulated beverage, included calcium, copper and zinc at a level of one-third the usual daily requirement. These are levels normally found in human diets.\nU.S. Pat. No. 3,992,555 issued to Kovacs (assigned Vitamins, Inc., 1976) describes food supplements prepared by mixing assimilable iron compounds, vitamins and minerals with a heated edible fat carrier. The minerals include calcium, zinc, copper, and manganese among others.\nU.S. Pat. No. 3,950,547 issued to Lamar et al (assigned Syntex Inc, 1976) describes a dietary composition containing peptides and/or amino acids, lipids and carbohydrates in an aqueous emulsion. Suitable minerals for adding at low levels include among others calcium, copper and zinc.\nU.S. Pat. No. 4,107,346 issued to Kravitz (1978) describes a dietary salt composition for use as a replacement for salt in foods. The role of copper is described as a component of several enzymes essential for nutrition. The patent further discloses that \"spontaneous fractures are common in animals feeding off copper deficient soils or who are given artificially depleted copper diets\". (column 4, lines 20-30) Zinc is described as helping in growth, wound healing and improving taste and smell. Clinical symptoms of manganese deficiency have not been observed in man. However, there is no question that manganese is essential for human nutrition. Main manifestations of its deficiency are impaired growth and skeletal abnormalities. Calcium is described as being necessary for blood coagulation, for calcium retention, and for relieving the symptoms of osteoporosis. Examples of salts that contain these four trace minerals are disclosed.\nU.S. Pat. No. 4,070,488 issued to Davis (unassigned, 1978) discloses a highly stabilized balanced nutritive composition useful in supplementing the diet of humans and/or animals. This composition contains gelatin. The patent discloses that the sulfhydryl groups of the gelatin can render copper inactive toward ascorbic acid.\nU.S. Pat. No. 4,214,996 issued to Buddemeyer et al (R.G.B. Laboratories, 1980) discloses mineral compositions which are very soluble. These compositions contain calcium, phosphorus, zinc, as well as manganese. Not all of the compositions that are described contain all four elements.\nU.S. Pat. No. 4,351,735 to Buddemeyer et al (R.G.B. Laboratories, 1982) is related to the '996 patent.\n\"Nutrients and Nutrition of Citrus Fruits,\" Citrus Nutrition and Quality, Ting, (American Chemical Society, 1980) discloses the presence of certain trace minerals in orange juice. These include copper, zinc, iron and manganese. Calcium and magnesium are the two major divalent cations in orange juice. The levels of all the minerals are low.\nHungerford et al, \"Interaction of pH and ascorbate in intestinal iron absorption,\" (1983) describes the iron absorption from various food materials. The diet which was low in iron also contained calcium carbonate, manganese sulfate and copper sulfate among others.\nU.S. Pat. No. 4,419,369 issued to Nichols et al (assigned Baylor College of Medicine, 1983) describes an improved dietary protein mineral module for infants. An approximate analysis of the material shows the presence of iron, zinc, copper and calcium.\nThe utility of these known supplements varies. Unlike agents (such as estrogen) which affect the metabolism of bone, calcium nutritional supplements have been thought to merely provide a source for calcium (which may or may not be properly absorbed and metabolized). See, for example, B. Riis et al., \"Does Calcium Supplementation Prevent Postmenopausal Bone Loss?,\" New England J. of Medicine, 316, 173-177 (1987); L. Nilas et al., \"Calcium Supplementation and Postmenopausal Bone Loss,\" British Medical Journal, 289, 1103-1106 (1984); and H. Spencer et al., \"NIH Concensus Conference: Osteoporosis,\" Journal of Nutrition, 116, 316-319 (1986).\nIt has now been discovered, however, that the administration of mixtures of certain calcium salts, along with trace minerals, copper, zinc and manganese are surprisingly effective for delaying age-related loss of bone. In particular, as compared to nutritional regimens known in the art, these methods afford greater efficacy in the treatment of age-related bone loss and related disorders.\nIt would be desirable, therefore, to have mixed calcium and mineral supplements which are compatible and nutritionally available. It would also be quite useful to have such supplements which could be added to food and beverage compositions without undesirably affecting organoleptic or aesthetic properties.\nIt is an object of the present invention to provide calcium mineral supplements which provide bone growth and can be used to treat age-related bone loss or to correct the imbalance that occurs between bone formation and bone resorption.\nIt is a further object of this invention to provide foodstuffs, beverages and beverage concentrates which are supplemented with calcium and trace minerals.\nThese and other objects are secured herein, as will be seen from the following disclosure."}
{"text": "FIG. 1 shows a front view of a parallel-flow-type condenser as an example of an aluminum alloy heat exchanger. Extruded multiple-holed tubes are used as tubes (extruded multiple-holed flat-shape tubes) 1, serving as refrigerant passages; each of the tubes 1 has a corrugated fin 2 brazed thereto, that is a radiator, and that is made of a brazing sheet fin clad with a filler alloy (parts of the repeated sections of the corrugated fins are omitted in FIG. 1); and the tubes 1 and the corrugated fins 2 are integrated by the brazing process, to produce a core 5, that forms the major part. 3 indicates header pipes, through which the tubes 1 communicate with one another. In recent years, in order to make heat exchangers small-sized and light in weight, some production methods are proposed. Among the methods is a production method in which, in place of the corrugate fins 2, which are limited in how thin they can be made, bear fins are used, which are thinner; the outside of the tubes 1 is covered with a filler alloy by metal-spraying, and they are combined and integrated by brazing (JP-A-63-34495 (\"JP-A\" means unexamined published Japanese patent application) and JP-A-2-84261, alternatively a method in which tubes and fins whose surfaces are metal-sprayed with a filler alloy are assembled into a core and brazing is carried out (JP-A-6-200344).\nWhen it is attempted to produce a heat exchanger by using tubes made by metal-spraying the above extruded tubes with a filler alloy, however, no heat exchanger is industrially satisfactory for practical use in brazing. The reason for that seems to be attributed to the occurrence of the following two problems with brazing.\nFirst is depletion of the filler alloy mentioned. This is a failure apt to occur when conventional brazing is carried out using brazing sheets having a smaller amount of a filler alloy, and this failure is liable to occur particularly in the case of tubes metal-sprayed with a filler alloy. Since the cost of metal-spraying with a filler alloy is high, when a coating of a filler alloy layer is formed by metal-spraying, the thickness of the filler alloy layer is generally 50 g/m.sup.2 or less, and this quantity is 1/2 or less of the quantity of the brazing sheet (the cladding ratio of the filler alloy being 10%), whose thickness is 0.4 mm, that is used generally as an electric-resistance weld pipe in heat exchangers. Accordingly, when the metal-sprayed filler alloy is not microscopically and uniformly deposited at the part where the amount of the filler alloy is small, depletion of the filler alloy sometimes occurs as a failure in which the part having a smaller amount of the filler alloy cannot be brazed (no fillet is formed).\nSecond is the occurrence of locally unbonded parts of heat exchangers. The occurrence of locally unbonded parts is a phenomenon peculiar to tubes metal-sprayed with a filler alloy, which is illustrated schematically in FIG. 2. When tubes 1 metal-sprayed with a filler alloy, and corrugated fins 2, are combined to form a core having several stages, the brazing can be carried out by heating, to effect the brazing with no problems. However, if a core 5 of multi-stage having about 40 stages is intended on a practical scale a part D, where the corrugated fin is not brazed to the tube at all, occurs at places. FIG. 4 indicates a side plate view.\nIf there is such a locally unbonded part D where a tube 1 and a fin 2 are not connected with each other, not only is the strength of the core 5 lowered as a whole, but also the sacrificial corrosion-preventive (cathodic-protection) effect of the tube 1 by the fin 2 cannot be secured, and the corrosion resistance is lowered greatly. The cause of the occurrence of the locally unbonded part D has not yet been fully clarified.\nThe most simple conceivable way for preventing the depletion of a filler alloy from occurring is a method wherein the amount of a filler alloy is increased satisfactorily. If there is a satisfactory amount of a filler alloy, the filler alloy goes to everywhere, even if the deposited amount of the filler alloy is uneven. When the amount of a filler alloy is increased, however, locally unbonded parts are liable to occur accordingly, as described below in the case of multi-stage one as mentioned above. Therefore, it is difficult to increase the amount of deposition of a filler alloy.\nTo solve this, it is conceivable to supply a filler alloy as uniformly as possible. The uniform deposition of a filler alloy should lead to the formation of fillets with a minimum amount of the filler alloy. However, the uniform deposition of a filler alloy was technically difficult in the conventional metal-spraying.\nTo prevent locally unbonded parts from occurring, it is recommended, in theory, to make the filler alloy applied thinly and to reduce microscopic projections and recesses (unevenness), so as to increase the filling density of the filler alloy. Namely, the thinner the filler alloy is, the smaller the change in the thickness of the tube before and after the brazing is made, and therefore locally unbonded parts are not likely to occur, even in a case of multi-stage. However, since thin application of a filler alloy means a decrease in the amount of the filler alloy, depletion of the filler takes place, even though locally unbonded parts can be prevented from occurring. Further, if the applied filler alloy has microscopic unevenness, gaps are liable to occur between the fin and the tube after the brazing, because the fin is in contact with the projections (projected areas) of the filler alloy when the core is assembled into a heat exchanger, as described above. The filling density of the filler alloy mentioned above is defined such that, when the filler alloy is present fully between the surface of a tube and the maximum height of the deposited filler alloy shown in FIG. 3 the filling density of the filler alloy is 100%. There is a tendency that the more deeply the unevenness of the filler alloy are formed, the smaller the filling density is. In the case of brazing sheets conventionally used in brazing, the filling density is about 100%. On the other hand, in the case of tubes metal-sprayed with a filler alloy, the metal-sprayed amount of the filler alloy is about 25 g/m.sup.2 (the possible minimum amount in brazing with a conventional Al/12% Si alloy), with the height of the filler alloy being 50 .mu.m or more (the height of the filler alloy (average filler alloy height) is about 8 .mu.m if the filler alloy is deposited with the filling density being 100%), and the filling density is 20% or less, which is a great decrease."}
{"text": "Due to the global climate crisis brought by the continuous increase of CO2 emission in the atmosphere (BP statistical Review of World Energy June 2016), the utilization of CO2 as a carbon resource is an important subject in the science and industry fields. The reduction of CO2 to utilizable C1 resources is a major research target in chemical conversions of solar energy, together with the hydrogen evolution by water reduction, to provide a strategic way responding the energy and environmental problems (I. Willner et al., J. Am. Chem. Soc., 109:6080-6, 1987). Among visible-light induced multi-electron reductions of CO2, the two-electron reduction to CO is a kinetically favourable choice because of the relatively low reaction barrier compared to one-electron and other higher-reduction reactions (A. J. V. Underwood, Ind. Eng. Chem., 32:449-54. 1940). Moreover, CO is currently utilized as a carbon source for the production of chemicals in petroleum chemical industry, particularly for the methanol production from a mixture of CO and H2, so-called syngas (M. E. Dry, Catal. Today, 71:227-41, 2002; I. Wender, Fuel Process. Technol., 48:189-297, 1996).\nSyngas has been also known as a key feedstock for the production of synthetic bulk chemicals via the Fischer-Tropsch (F-T) processing (A. J. V. Underwood, Ind. Eng. Chem., 32:449-54. 1940). For the production of chemicals from syngas, the tailoring of H2/CO ratio in syngas mixtures is critical, e.g. 2:1 H2/CO for methanol production and F-T hydrocarbon syntheses (M. E. Dry, Catal. Today, 71:227-41, 2002; K. C. Waugh, Catal. Today, 15:51-75, 1992) and 1:1 H2/CO for the production of aldehydes via hydroformylation of alkenes (M. Beller et al., J. Mol. Catal. A: Chem., 104:17-85, 1995). While the simultaneous formation of H2 and CO is known to occur by photochemical and electrochemical reductions of CO2 and water using Ni— (V. S. Thoi et al., J. Am. Chem. Soc., 135:14413-24, 2013), Re— (B. Kumar et al., Chem. Commun., 45:272-4, 2012), and Ru— (P. Kang et al., Energy Environ. Sci., 7:4007-12, 2014)-molecular catalysts, little has been referred to possible control of syngas compositions. Therefore, a new protocol that focuses on adjustability of H2/CO ratio in visible-light induced syngas production should be of scientific significance.\nConventional technologies for catalytic production of syngas require the production of hydrogen and carbon monoxide through independent catalytic conversion processes and mixing of the gases through additional processes. Further, most of the production technologies are limited to the utilization of by-product gases and related processes thereof. Recent technologies for direct production of syngas using molecular catalysts are based on electrochemical methods that still need external energy.\nThus, the present inventors have earnestly and intensively conducted research to develop a method for direct production of syngas without the need for external energy and a catalytic system for producing syngas in which the ratio of gases formed can be controlled, and as a result, found that when two types of transition-metal complex catalysts, i.e. a rhenium catalyst and a cobalt catalyst, are immobilized on TiO2 particles, the resulting catalytic system exhibits high catalytic activity, produces H2 and CO without external energy, and enables control over the ratio of H2/CO formed. The present invention has been accomplished based on this finding."}
{"text": "The present invention pertains generally to target scoring devices and more particularly to a device for communicating an intercept score in a noisy environment. In the testing of both equipment and human skill, it is oftentimes desirable to automatically signal that a shell or projectile has come within the zone of the maximum detection range of a miss distance indicator of a target which is being fired upon. Penetration of the zone is commonly referred to as a score. This situation often occurs in ship gunnery exercises where a ship's crew wishes to determine the number of shells which have penetrated this zone.\nConventional automatic scoring systems have used various methods for detecting scores. A common system has used a doppler detector in which alternating current cycles amounting to heat frequency oscillations are generated by a miss distance indicator and are used to modulate a telemetry transmitter. The telemetry signal is, in turn, received by a receiving station on board the ship and demodulated so that the beat frequency oscillations can be detected. The oscillations generally range in a frequency from a few Hertz up to as much as 25 kiloHertz. The oscillations are then fed through a detection system which detects the presence of the oscillations above a certain threshold frequency. A counter is then used to count these oscillations above this threshold frequency to discriminate against noise and thereafter register a score.\nIn practice this system has presented various disadvantages and limitations in the presence of environmental and radio frequency noise. For example, in a case where the receiving antenna of a scoring indication unit is located proximate various radar transmitting antennas and other forms of RF interference, false score indications are prevalent. In this case and in the loss of the telemetry carrier, the automatic scoring device at the receiving station cannot effectively distinguish between the desired signal and the noise. The system, in general, has therefore not shown the reliability required to properly evaluate the functions being tested."}
{"text": "This invention relates generally to surgical wound retractors and more specifically to circular wound enactors that provide circumferential retraction of an incision or wound and also provide isolation of the surgical incision or wound from the lumen of the retractor.\nRetraction and isolation of a surgical incision can be an important element associated with a surgical procedure. Adequate access may be provided by a circumferential retraction device so that a surgeon may operate in a clear and open, generally unrestricted field. In addition, a circumferential retractor having a fluid impermeable sleeve extending through the incision or wound may prevent excessive bleeding or contamination of adjacent tissue. In some surgeries, it is important to avoid contaminating the incision or wound site.\nCommon procedures for placing current circumferential wound retraction devices involve placement of an internal retention ring that is coupled to a thin-walled sleeve or tube that is further coupled to an external retention member. The external retention member is wound or turned upon itself to wind the sleeve upon the external retention member. The physics associated with winding a sleeve upon a circular winding collar or external retention member dictate that ample energy must be applied to the external winding motion to provide adequate circumferential retraction of a linear incision. The initial windings of the process are not overwhelming, however, as the process continues the retraction becomes more difficult and the mechanical advantages first appreciated begin to dissipate. The final windings are more difficult since a balance between the force required to wind the sleeve upon the external retention member must be balanced with the force acting to unwind the sleeve from the external retention member.\nThere remains a need to provide a substantially rigid, noncompliant circumferential retractor that easily shortens the length of an associated tubular sleeve coupled to an internal retention member so that the shortening of the sleeve results in dramatic retraction of an incision or wound. Additionally, the retention members must remain substantially circular so that the incision is retracted symmetrically."}
{"text": "The physical properties of polyethylene terephthalate are such that this material lends itself well to the injection stretch blow molding process.\nIn contrast, the physical properties of polyethylene are considered much less well-suited for injection stretch blow molding. JP-A-2000/086722, published on Mar. 28, 2000, discloses a high density polyethylene resin which is subjected to injection stretch blow molding. The polyethylene resin has a melt flow index of from 1 to 15 g/10 minutes; a flow ratio of from 10 to 14.5; and a density of from 0.961 to 0.973 g/cm3.\nPlastic parts stress crack when they are under tensile stress and in contact with liquids either by an oxidative or lubrication mechanism. This process is known as Environmental Stress Cracking. The oxidative mechanism (i.e., cleavage of polymer molecules) is present in liquids that contain oxidants (e.g. hypochlorite, and hydrogen peroxide) and the lubrication mechanism is present in liquids that contain surfactants. With either mechanism, the number density and length of the tie molecules (i.e., molecules that connect the various crystallites) is the controlling parameter for the environmental stress crack resistance (ESCR). Note that stress cracking occurs only in tensile loading, not compressive. In other words, in a bottle, stress cracking can occur only in the regions that are under tensile deformation and in contact with the liquid. As the polyethylene material is under tensile load in the trade or in the lab, the crystallites are under stress and they start moving away from each other as the tie molecules are stretched. In the brittle mode of failure, the tie molecules completely pull away from the crystallites causing them to separate. In the ductile failure, the tie molecules cause the crystallites to fracture and create fibrils. The oxidants in the liquid (e.g. bleach) cleave the tie molecules causing earlier failure than when the material is exposed to water or air. Furthermore, surfactants in the liquid lubricate the disentanglement of the tie molecules and their separation from the crystallites. In terms of measurable properties, ESCR increases with the average molecular weight (as the number of tie molecules increases with molecular weight), and decreases with increasing breadth of the molecular weight distribution, crystallinity and spherulitic size.\nLimitations including manufacturing cost and manufacturing speed, and environmental stress crack resistance properties have prevented injection stretch blow molded polyethylene bottles from achieving commercial success. The present invention aims to overcome the limitations of the prior art."}
{"text": "In wireless communication systems, one base station (BS) generally provides services to a plurality of user equipments (UEs). The BS schedules user data for the plurality of UEs, and transmits the user data together with control information containing scheduling information for the user data. In general, a channel for carrying the control information is referred to as a control channel, and a channel for carrying the user data is referred to as a data channel. The UE finds control information of the UE by searching for the control channel, and processes data of the UE by using the control information.\nIn order for the UE to receive user data assigned to the UE, control information for the user data on a control channel must be received. In a given bandwidth, a plurality of pieces of control information for a plurality of UEs are generally multiplexed within one transmission interval. That is, to provide a service to the plurality of UEs, the BS multiplexes the plurality of pieces of control information for the plurality of UEs and then transmits the control information through a plurality of control channels. The UE searches for control channel of the UE among the plurality of control channels.\nBlind decoding is one of schemes for detecting specific control information from the plurality of pieces of multiplexed control information. The blind decoding attempts to recover a control channel by using several combinations of information in a state where a UE has no information required to recover the control channel. That is, in a state where the UE does not know whether control information transmitted from the BS is control information of the UE and the UE does not know in which portion the control information of the UE exists, the UE decodes all pieces of given control information until the control information of the UE is found. The UE can use information unique to each UE to detect the control information of the UE. For example, when the BS multiplexes control information of each UE, an identifier unique to each UE can be transmitted by being masked onto a cyclic redundancy check (CRC). The CRC is a code used for error detection. The UE de-masks unique identifier of the UE from the CRC of the received control information, and then can detect the control information of the UE by performing CRC checking.\nIf the UE cannot correctly detect the control information of the UE from the pieces of multiplexed control information, user data on the data channel cannot be decoded. Therefore, fast and correct detection of the control information has a significant effect on overall system performance. However, it may be difficult to detect the control information when using only blind decoding. Since each UE may require different control information and may use a channel encoding scheme by the use of a different code rate, each UE may have control information with a different size. Therefore, the number of blind decoding attempts may be significantly increased in a control region in which control information is transmitted. Battery consumption of the UE increases in proportion to the number of detection attempts.\nMeanwhile, as a mobile communication system of a next generation (i.e., post-3rd generation), an international mobile telecommunication-advanced (IMT-A) system is standardized aiming at support of an Internet protocol (IP)-based seamless multimedia service in an international telecommunication union (ITU) by providing a high-speed data rate of 1 gigabits per second (Gbps) in downlink communication and 500 megabits per second (Mbps) in uplink communication. In a 3rd generation partnership project (3GPP), a 3GPP long term evolution-advanced (LTE-A) system is considered as a candidate technique for the IMT-A system. The LTE-A system is evolved to increase a completion level of the LTE system, and is expected to maintain backward compatibility with the LTE system. This is because the provisioning of compatibility between the LTE-A system and the LTE system is advantageous in terms of user convenience, and is also advantageous for a service provider since existing equipment can be reused.\nAccordingly, there is a need for a method in which a BS effectively transmits control information and a UE effectively receives data by using the control information in a next generation mobile communication system such as an LTE-A."}
{"text": "Eye diagrams are a conventional format for representing parametric information about signals, and especially digital signals. We shall refer to an item of test equipment or a measurement circuit arrangement that creates an eye diagram as an eye diagram tester, whether it is found in an oscilloscope, a BERT (Bit Error Rate Tester), a logic analyzer, or, as a separate item of test equipment. The preferred method (and by implication, any corresponding circuit apparatus) to be disclosed herein is that of the incorporated Application, is different than those used in \"\"scopes and BERTs, and is especially suitable for use within a logic analyzer, as well as an item of stand-alone test equipment. We shall call this different method (and any corresponding circuit apparatus) an Eye Diagram Analyzer, or EDA for short.\nA modern eye diagram for a digital signal is not so much a trace formed continuously in the time domain, as it is an xe2x80x9ceyexe2x80x9d shape composed of closely spaced points (displayed dots, or illuminated pixels) representing many individual measurement samples taken upon separate instances of a signal occurring on a channel of interest, and which were then stored in a memory. Each measurement sample contributes to a displayed dot. The eye shape appears continuous because the collection of dots is rather dense, owing to the large number of times that the signal is sampled. Unlike a true continuous technique, however, there may be detached dots that are separated from the main body of the eye shape. A further difference with the continuous analog technique is that rare or infrequently occurring events, once sampled, do not appear faint in the display or disappear with the persistence of the CRT\"\"s phosphor. This latter difference is often quite an advantage, since it is often the case that such otherwise xe2x80x9chard to seexe2x80x9d features of the trace are very much of interest.\nIn any event, the vertical axis is voltage, and the horizontal axis represents the differences in time (i.e., various offsets) between some reference event and the locations for the measurement samples. The reference event is generally an edge of a clock signal in the system under test, represents directly or through some fixed delay the expected point in time when the value of an applied data signal would be captured by some receiving circuit in an SUT (System Under Test), and is derived from an application of the SUT\"\"s clock to the Eye Diagram Analyzer. The time axis will generally have enough length to depict one complete eye-shape (cycle of a DUT signal) centered about the reference, with sometimes perhaps several additional eyes (cycles) before and after.\nDifferent (X, Y) regions of an eye diagram represent different combinations of time and voltage. Assume that the eye diagram is composed of a number of pixels, and temporarily assume that the resolution is such that each different (X, Y) pixel position can represent a different combination of time and voltage (and vice versa), which combinations of time and voltage we shall term xe2x80x9cmeasurement points.xe2x80x9d What the analyzer measures is the number of times, out of a counted number of clock cycles, that the signal on the channel being monitored passed through a selected measurement point. Then another measurement point is selected, and the process repeated until there are enough measurement points for all the pixels needed for the display. The range over which the measurement points are varied is called a xe2x80x9csample spacexe2x80x9d and is defined during a measurement set-up operation. And in reality, we define the sample space and the resolution for neighboring measurement points first, start the measurement and then let the analyzer figure out later how to ascribe values to the pixels of the display. The xe2x80x9cdisplayxe2x80x9d is, of course, an arbitrary graphic output device such as a printer or an X Window of some as yet unknown size in a window manager (e.g., X11) for a computer operating system. (A one-to-one correspondence between display pixels and measurement points is not required. It will be appreciated that it is conventional for display systems, such as X Windows, to figure out how to ascribe values to the pixels for an image when the correspondence between the display\"\"s pixel locations and the image description is not one-to-one.)\nAn eye diagram trace itself is not a single time domain waveform (think: xe2x80x98single valued functionxe2x80x99), but is instead equivalent to an accumulation of many such instances; it can present multiple voltage (Y axis) values for a given time value (X axis). So, for example, the upper left-hand region of an eye might represent the combination of an adequate logical one at an adequately early time relative to the SUT\"\"s clock signal, and an eye diagram whose trace passes robustly through that region indicates to us that a signal of interest is generally achieving a proper onset of voltage at a proper time. Furthermore, we note that there are also other regions, say, near the center of an eye, that are not ordinarily transited by the trace, and which if that were indeed to happen, would presumably be an indication of trouble. Thickening of the traces is indicative of jitter, a rounding of a corner is indicative of slow transitions, and so on. An eye diagram by itself cannot reveal in the time domain which isolated instance of the signal caused such an exception, as other types of measurements might, but it does provide timely and valid information about signal integrity within a system as it operates. In particular, by incorporating very long (perhaps xe2x80x9cinfinitexe2x80x9d) persistence the eye diagram presents readily seen evidence of occasional or infrequently occurring failures.\nAn eye diagram, then, is information about signal behavior over time at various time-voltage (X, Y) combinations. A simple system would be to indicate that the signal was xe2x80x9ctherexe2x80x9d or that it wasn\"\"t. That is, respectively put either an illuminated pixel or a non-illuminated pixel at the various (X, Y) locations for the different instances of xe2x80x9cthere.xe2x80x9d This is similar to what an analog oscilloscope would do if it were used to create an eye diagram for some signal. However, in such an analog case we would notice that some parts of the trace were brighter than others, and understand that this is a (somewhat useful) artifact caused by finite persistence on the one hand (old stuff goes away) and relative rates of occurrence on the other. That is, the display ends up having an intensity component at each pixel location. This is fine as far as it goes, but we would rather not rely upon the persistence of phosphor for this effect, since the most interesting indications are apt to be also the faintest. Since we are not using an analog \"\"scope, anyway, and have an instrument (an EDA) with memory (akin to a digital oscilloscope, timing analyzer or logic analyzer), we can gather data and decide after the fact what pixel value is to go with each (X, Y) pixel location. Those pixel values can be rendered as variations in color, intensity, or both, according to whatever scheme is in use (and there are several). The general idea is that the operator configures the EDA to render the display in a way that makes the condition he is most interested is quite visible, and also such that the eye diagram as a whole is generally easy to interpret. Thus, the reader is reminded that there is usually more going on than simply the presence or absence of dots at some series of (X, Y) locations, and that we often have recourse to the notion of a xe2x80x9cpixel valuexe2x80x9d at some (X, Y) position in the display. We shall denote with the symbol xe2x80x9cVxe2x80x9d whatever that xe2x80x9cpixel valuexe2x80x9d is. V might be a single binary-valued item, such as xe2x80x9cONxe2x80x9d or xe2x80x9cOFFxe2x80x9d or it might be a variable intensity without the notion of color. On the other hand, V will often expand into a whole series of parameters VR, VG, VB, . . . , where each such parameter represents the intensity of an associated color. Whatever V is or represents, each Vi is computed by the Eye Diagram Analyzer according to operational modes selected by the operator and in response to events measured in the System Under Test. In any event, we shall simply write (X, Y, V)i or perhaps (Xi, Yi, Vi), depending upon which notation seems to work best for our needs at the time, and not bother to indicate any associated further expansion of a Vi into its component parameters, such as (VR, VG, VB).\nIt is often the case that the utility of an eye diagram is needed for characterizing or discovering circumstances that are both erroneous and very occasional. It is also the case that some SUTs have a great many channels that are subject to investigation. Some busses have hundreds of member signals, for example. When faced with such circumstances, the xe2x80x9cluxuryxe2x80x9d of having one individual eye diagram trace per SUT signal becomes an oppressive burden. We might measure it that way, and we can indeed display it that way (with four or maybe eight channels at a time), but we likely will have lost all patience and become disgusted with the whole process before we complete looking at twenty-five or more sets of four or more traces each. Surely that is the wrong way to go about analyzing the data! But on the other hand, automating the investigation may be risky. Masking measurements, for example, essentially require that we formally decide ahead of time what is not of interest. The analyzer can apply the mask for us automatically and at great speed, but we will never know for sure that there was not some irregularity in there that met the mask criteria, and yet that would have been of interest to us anyway, if we had only seen it.\nAccordingly, another tool has emerged to assist in eye diagram analysis for situations involving many channels. It is to merge into one combined eye diagram presentation the data of separate xe2x80x9ccomponentxe2x80x9d eye diagram measurements for a useful grouping of related signals. The operator can instruct the EDA to merge the data of individual eye diagrams for a specified plurality (family) of signals into one combined eye diagram by arithmetically combining the numerical values for the sampled component data. Note that this does not involve any notion of adjusting vertical position for separately displayed component eye diagrams to get them to xe2x80x9cstack upxe2x80x9d! So now we have a combined eye diagram that probably has fat traces (indicating that, as expected, not all signals have simultaneous and identical rise times, voltage levels, etc.). Independent of that, we now expect that, having merged everything together, if there is something unusual going on, even if only occasionally for just one channel, we will, in principle, be able to see it.\nWe shall term such a combined eye diagram, created from the merging of data for individual component eye diagrams, a xe2x80x9ccompositexe2x80x9d eye diagram.\nThe particular eye diagram analyzer described in the incorporated Application xe2x80x9cCOMPOSITE EYE DIAGRAMSxe2x80x9d can assign a plurality of SUT data signals to be members of a labeled group of channels. There may be a plurality of such groups. In addition to mere superposition in an (X, Y) display space of the various i-many (X, Y)-valued pixels for individual component eye diagrams associated with that group, other measured data for those pixels within a group can be merged in any of various different modes to produce corresponding composite eye diagram presentations. In general, the rule for combining components to obtain a composite will produce a series of (X, Y, V)i, where the Vi represent some sort of density. The Vi are generally rendered with variations in intensity or color, or with some combination of both. The details of that would be tedious for achieving our purpose in this Application, and interested readers are referred to the incorporated xe2x80x9cCOMPOSITE EYE DIAGRAMSxe2x80x9d.\nThis merging of measurements in the displayed data of a composite eye diagram is fine as far as it goes. It greatly assists in discovering that there is some trouble in there, someplace. It allows us to cast a wide net, as it were. But the very property that give us confidence that xe2x80x9cit\"\"s in there, someplacexe2x80x9d also poses the dilemma of how to separately identify it when its xe2x80x9chiding in plain sight,xe2x80x9d so to speak. Say, at least one signal out of thirty-two was unable to drive its load. Presumably, there will be a visible indication of this situation in the shape of the composite eye diagram. But on account of which channel(s)?? We need a way to quickly and easily isolate the culprit and discover who it is. Fussing with the probes one signal at a time is almost never practical, and in some cases, is impossible (as when an entire bus is probed at once by a single multi-conductor connector).\nOnce we have a channel in mind (whether for a component or a composite), and are interested in a particular place within an eye diagram, it would also be useful if there were a way to discover what the relevant history of signal behavior for that channel is at that place in the eye diagram. That is, it would be desirable to discover the answers to such questions as: xe2x80x9cHow many times did it contribute to that part of the eye diagram? What are the exact times and voltages for that location? How many clock periods occurred during this measurement?xe2x80x9d\nEvidently, then, we are in need of some tools to assist in the understanding and analysis of eye diagrams, and especially of those that are presented as composites. What to do?\nA composite eye diagram is one where the data for a group of several selected individual component eye diagram measurements are combined to produce a single presentation, whose displayed pixels typically represent some sort of density. In a system that produces and displays composite eye diagrams, the operator may position a screen pointer over some interesting part of the composite eye diagram, and the analyzer will display a list of those channels that contributed to the displayed pixels in that part of the eye diagram, as well as information in tabular form related to the behavior of those channels at that part of the eye diagram. The operator may subsequently, or at any other time, select a channel from among the group, and then have the displayed pixels for which that channel is responsible be displayed in a selected color or otherwise highlighted in a way that allows them to be distinguished from all others. So, for an unusual indication in a composite eye diagram the user can discover that channels 2, 5 and 8 are involved by pointing at that indication with the mouse pointer, and then by highlighting the individual traces for those channels, decide that activity for channels 2 and 8 is benign, but that the activity for channel 5 is non-conforming."}
{"text": "Omnidirectional antennas are used for example as indoor antennas. They are multiband capable and preferably radiate with a vertical polarisation orientation. For this purpose, they may comprise a base or earth plate (reflector), which may for example be disc-shaped and on which a monopole radiator rises transversely and in particular perpendicularly to the base plate. The entire arrangement is generally covered by a protective housing, i.e. an antenna cover (radome).\nThe present broadband omnidirectional antenna can not only be used within buildings, but for example also in vehicles, in particular rail vehicles or boats.\nA generic omnidirectional antenna is known for example from DE 103 59 605 A1. The monopole radiator known from this document rises vertically above a base plate, from which it is galvanically isolated. The antenna known from this document comprises a vertically polarised monopole radiator. In this case, the vertically polarised radiator is in particular in the shape of a hollow cylinder or hollow cone and extends away from the base plate.\nThe omnidirectional antenna from DE 103 59 605 A1 is disadvantageous in that the lower limiting frequency is limited by the specified overall height and the specified diameter.\nThe example non-limiting technology provides a broadband omnidirectional antenna which can be produced so as to be as simple, cost-effective and compact as possible, and which at the same time covers a wider frequency spectrum.\nThis is achieved by means of a broadband omnidirectional antenna as described herein.\nA broadband omnidirectional antenna comprises a first radiator that is arranged on a base plate, which base plate is preferably also used as a reflector, and that has a longitudinal axis which extends at least approximately, predominantly or completely perpendicularly to the base plate. In that case, the first radiator extends from the base plate away therefrom. The first radiator has a first end comprising a foot and/or feed-in point and a second end which is opposite the first end. The first end, i.e. the foot and/or feed-in point, of the first radiator is in this case galvanically isolated from the base plate, but is arranged closer to the base plate than the second end. The first radiator also comprises radiator surfaces which originate in the region of the first end and extend towards the second end. A distance between the radiator surfaces and the longitudinal axis increases at least in portions from the first end towards the second end. This means that the radiator surfaces diverge from one another along the longitudinal axis at least over a partial length. Furthermore, the omnidirectional antenna comprises a second radiator which comprises at least one radiator surface. The second radiator is arranged on the first radiator so as to be galvanically isolated therefrom and can be fed preferably exclusively or predominantly by the first radiator. In one embodiment, the radiator surfaces of the second radiator are arranged in relation to the radiator surfaces of the first radiator such that they can act as a continuation thereof. This means that the second radiator is a continuation of the first radiator. In this case, the radiator surfaces of the second radiator can be inclined at least in portions or can only extend in parallel with the longitudinal axis. They are spaced further apart from the base plate than the radiator surfaces of the first radiator. Alternatively, i.e. in another embodiment, it would also be possible for the at least one radiator surface of the second radiator to be arranged in the region of the second end of the first radiator, in particular between the radiator surfaces of the first radiator, i.e. within said radiator, so as to be in parallel with the base plate or such that one of the components thereof is predominantly parallel to said base plate.\nIt is particularly advantageous for the second radiator to be fed exclusively or predominantly by the first radiator. In this case, a separate feed line for the second radiator is not required or provided. In this case, it is advantageous for the second radiator to be a continuation of the first radiator, the two radiators being galvanically isolated from one another. This increases the band width that can be produced and keeps the production costs low.\nIn an advantageous embodiment of the broadband omnidirectional antenna, a feed device is arranged at the foot and/or feed-in point. In this case, the feed device extends towards the base plate and preferably passes therethrough. A connector element, in particular in the form of a socket, is arranged on a bottom side of the base plate, which side is opposite the assembly side comprising the received first and second radiators. A feed cable can be or is connected to said connector element. The feed device preferably extends, at least by its first end, into the connector element, it being possible for electrical contact to be established, or said electrical contact being established, at least indirectly (via an additional conductor) or directly, between the first end of the feed device and an internal conductor of the feed cable. In this case, the feed device is galvanically isolated from the base plate. Depending on the embodiment of the broadband omnidirectional antenna, the feed device is galvanically, but preferably in a solder-free manner, connected to the first radiator at the foot and/or feed-in point. The feed device could also be capacitively coupled to the first radiator at the foot and/or feed-in point, the feed device extending towards the second end of the radiator surfaces of the first radiator at least in part along the longitudinal axis or such that one of its components is predominantly in parallel with the longitudinal axis.\nIn this case, it is particularly advantageous for the foot and/or feed-in point of the first radiator to have a sleeve-shaped or hollow cylindrical extension towards the second end of the first radiator. The feed device is arranged in the sleeve-shaped extension at least over a partial length thereof, the feed device and the sleeve-shaped extension being galvanically isolated from one another. The sleeve-shaped extension can extend as far as the second end of the first radiator or beyond the second end of the first radiator. Depending on the use, the first radiator can thus be fed capacitively or inductively.\nIn a particularly preferred embodiment, the first radiator has, along its longitudinal axis and over its entire length or a partial length thereof, a progression that is in part or predominantly or completely conical or funnel-shaped. The second radiator comprises a predominantly or preferably completely peripheral radiator surface, a diameter or circumference of the peripheral radiator surface of the second radiator at the first end thereof being adapted to a diameter or circumference of the second end of the first radiator.\nAdaptation of this kind is preferably achieved by the diameter or circumference at the first end of the second radiator deviating from the diameter or circumference at the second end of the first radiator by less than 20%, 15%, 10%, 8%, 5% or 3%. It is particularly advantageous for the diameter or circumference at the first end of the second radiator to be slightly larger than the diameter or circumference at the second end of the first radiator. “Slightly larger” should be understood to mean larger by a small number of millimeters, in particular by less than 8 mm, 6 mm, 4 mm or 2 mm, but preferably by more than 1 mm, 3 mm, 5 mm, 7 mm or 9 mm.\nIn the context of another embodiment, the diameter of the second radiator remains constant along the longitudinal axis or decreases in the direction of the longitudinal axis from the first end towards the second end. This is particularly advantageous in that the omnidirectional antenna can be constructed so as to be compact.\nIn another preferred embodiment of the omnidirectional antenna, the second radiator comprises one or more slots, which extend from the second end thereof, which is opposite the first end, towards said first end and terminate at a distance therefrom. In this case, the width of these slots can be constant or decrease towards the first end. In principle, the first slots could also extend from the first end towards the second end and terminate at a distance from the second end.\nSo that the first radiator and the second radiator are permanently oriented relative to one another in a precisely defined position, in a particularly preferred embodiment of the omnidirectional antenna, a (dielectric) holding and/or spacing element is used which is arranged at least in part within the first radiator and is non-rotatably fastened thereto. The holding and/or spacing element is preferably also non-rotatably fastened to the second radiator, the holding and/or spacing element being designed such that a gap (along the longitudinal axis) between the first end of the second radiator and the second end of the first radiator has an adjustable width. The first radiator and the second radiator are therefore arranged in relation to one another such that they do not overlap. The holding and/or spacing element therefore performs a number of functions. Firstly, the holding and/or spacing element prevents the first radiator and the second radiator from rotating relative to one another over time. Furthermore, said element ensures that the first radiator and the second radiator are galvanically isolated from one another. The gap, which is adjusted between the first radiator and the second radiator by the holding and/or spacing element, is preferably larger than 0.1 mm, 0.3 mm, 0.5 mm, 0.7 mm, 0.9 mm, 12 mm, 15 mm, 17 mm, 20 mm, 30 mm, 40 mm or 50 mm, and is preferably smaller than 40 mm, 30 mm, 20 mm, 18 mm, 16 mm, 13 mm, 11 mm, 9 mm, 8 mm, 6 mm, 3 mm or 1 mm.\nIn another preferred embodiment, the first radiator comprises n radiator surfaces, where n≥2. In this case, the n radiator surfaces are galvanically interconnected or formed in one piece with one another at the first end of the first radiator, the radiator surfaces being arranged around the longitudinal axis of the first radiator so as to be offset from one another, thus forming slots between adjacent radiator surfaces, and the slots beginning at a distance from the first end of the first radiator and extending as far as the second end of the first radiator. In this case, at least part of the at least one radiator surface of the second radiator is arranged at the second end of the first radiator, between the radiator surfaces of the first radiator, so as to be in parallel with the base plate or such that one of the components thereof is predominantly in parallel with said plate. What is particularly advantageous here is that a radiator arrangement of this kind can be produced in a very simple manner, for example from sheet metal parts. An omnidirectional antenna of this kind has a very low overall height, but still operates at a wide range of frequencies.\nIn another embodiment, the radiator surfaces of the first radiator comprise a plurality of radiator partial surfaces which are oriented at an angle to one another. The same can also apply to the at least one radiator surface of the second radiator.\nIn this case, the radiator surfaces of the first radiator and second radiator are preferably free of curves (except for the bending edge) and are each arranged in a separate plane. In this case, the first radiator and/or the second radiator can be produced from a metal sheet in a cutting, stamping and/or bending process.\nIn a particularly preferred embodiment of the omnidirectional antenna, said antenna comprises a coupling device. The coupling device is used in order for it to be possible for the lower limiting frequency at which the omnidirectional antenna can be operated to be reduced further. For this purpose, the coupling device comprises one or more coupling projections, a first end of the coupling projection or coupling projections being galvanically connected to the radiator surface of the second radiator and extending towards the base plate. The coupling projection or coupling projections is/are spaced further apart from the longitudinal axis than the radiator surfaces of the first radiator and second radiator. This means that the coupling projection or coupling projections extend towards the base plate outside of the first radiator and second radiator. At least one coupling surface is formed or integrally formed on a second end of the coupling projection or coupling projections that is opposite the first end and is therefore arranged closer to the base plate than said first end, which coupling surface is galvanically connected to the relevant coupling projection. The at least one coupling surface extends in parallel with the base plate or such that one of the components thereof is (predominantly) in parallel with said plate. Owing to coupling of this kind that is relative to the base plate, the lower limiting frequency can be reduced further. In this case, it is possible for the omnidirectional antenna to be operated in a frequency range of 600 MHz to 6 GHz. Said antenna is preferably operated in a frequency range of 650 MHz or 698 MHz to 6 GHz. Depending on the size and dimensions of the feed point, inter alia, it is also possible for the frequency range to be widened at the upper and/or lower limit.\nIn another embodiment of the omnidirectional antenna, the at least one coupling surface is galvanically connected to the base plate or is arranged at a distance therefrom such that the at least one coupling surface is capacitively coupled to the base plate. The distance between the coupling surface and the base plate and the size of the coupling surface can be varied as desired, depending on the use. The coupling surface can be arranged so as to be in parallel with the base plate. It can also be arranged obliquely or designed so as to be uneven (e.g. undulating).\nIn this case, an additional dielectric can be arranged between the at least one coupling surface and the base plate, for example, on which dielectric the at least one coupling surface rests or is supported. As a result, the coupling can again be adjusted more accurately and the stability of the omnidirectional antenna as a whole can be increased.\nIn another preferred embodiment, the plurality of coupling projections are galvanically connected to a common coupling surface by means of the second end thereof, the coupling surface being in the form of a common coupling frame which defines a receiving space in which part of the first radiator is arranged. In principle, the common coupling frame can be of any shape. In particular, the cross section thereof may be rectangular, square, circular or oval.\nIn order to further increase the stability of the omnidirectional antenna and further increase weather resistance, in another embodiment, said antenna comprises a covering hood. Preferably, one single covering hood is used, which is connected to the base plate in an interlocking and/or frictional and optionally moisture-tight manner, and surrounds the first radiator and second radiator. In this case, the covering hood is preferably arranged such that it is not in contact with the first radiator and the second radiator."}
{"text": "1. Field of Invention\nThe present invention relates to a printing device that prints on both sides of a print medium and has a page mode for converting one page of print data after setting front and back side print data conversion areas in a data storage device, and to a control method for the printing device.\n2. Description of Related Art\nPrinting devices (simply “printers” below) for printing on both sides of paper or other print medium (commonly called “duplex printing”) are taught in Japanese Unexamined Patent Appl. Pub. JP-A-03-027434 and Japanese Unexamined Patent Appl. Pub. JP-A-04-001067.\nJapanese Unexamined Patent Appl. Pub. JP-A-03-027434 and Japanese Unexamined Patent Appl. Pub. JP-A-04-001067 teach a printing system that has a repagination unit, an odd-page processing unit, an even-page processing unit, and a reverse-side print processing unit. The repagination unit repaginates and updates the page numbers when a duplex printing command is issued by the host computer. The odd/even page determination unit determines whether a page number is odd or even. If the page number is odd, the odd-page processing unit prints on the front of a sheet and then discharges the sheet into an intermediate holding tray. If the page number is even, the even-page processing unit writes the start address and length of the even page print command string to a control table. The reverse-side print processing unit prints the even pages on the second side of the sheets on which the odd pages have been printed on the first side and are stored in the holding tray.\nThe host computer issues a duplex printing command to the printing system in order to print on both sides. More particularly, a command must be asserted from the host computer in order to switch between simplex printing and duplex printing. Printers that switch between duplex printing and simplex printing according to the on or off state of a DIP switch are also known from the literature. If the storage capacity of the intermediate holding tray is exceeded during the first-side printing process or the end of a single-sided document is detected when the duplex printing mode is active, operation switches from the first-side printing process to the second-side printing process.\nPrinters that print to continuous paper such as roll paper are also known. When such printers are used for printing sales receipts to roll paper, the length of the printed data is not constant and the length of each issued receipt varies with each sale transaction. The issued receipt is long when there are many purchased products and there is much information to print, and is short when there are few products (little information). If the roll paper printer can print on both front and back sides of the roll paper and can print long receipts in the duplex printing mode, the length of the issued receipt can be shortened and roll paper consumption can be reduced.\nOne type of printer used for printing receipts is the line thermal printer. In general, each time a line thermal printer receives print data including one line of print data and a carriage return, it prints the one line and then advances the paper one line. If a carriage return is not included but print data exceeding one line is received, the printer also executes the printer process and advances the paper. One page (one receipt) is printed and output by repeating these steps. This printing process mode is referred to herein as the “standard mode.”\nDemand for printing receipts and coupons having a greater visual impact on the customer has also grown in recent years. Printers with a wide range of print functions, such as for printing logos, images, and barcodes in addition to text data, and changing the output format by changing the print orientation (page orientation), are also common. In order to provide such diverse printing functions, line printers that have a page mode enabling printing a free-form layout in addition to the standard mode described above are also known from the literature. When the page mode is selected, a prescribed printing area is set in a buffer (memory) enabling converting print data (that is, conversion to data enabling printing by the print head) for at least one full page, all of the print data is converted into the set printing area, and a start printing command is then asserted to start a batch printing process to print all of the data in a single process (print the converted print data by the print head).\nThus setting the printing area enables freely arranging the layout of the print data in the set printing area during conversion. For example, a horizontal layout in which the normal print orientation is horizontal (across the width of the page) can be converted to a vertical layout in which content is printed in the paper feed direction by rotating the print data 90 degrees in the print buffer when converting the print data. The page mode thus enables print functions such as page layouts that cannot be printed in the standard mode, and thus enables more effective receipt printing.\nWhen a printer capable of duplex printing prints both sides, it is also desirable to shorten the length of the issued receipt and reduce paper consumption even when operating in the page mode. In duplex printing, however, the print result is affected by how the first-side data and the second-side data are allocated to the front and back sides. For example, if the break between the first-side data and the second-side data occurs in the middle of an image, the image will be broken into two parts and a receipt or coupon that has a positive visual impact on the user cannot be printed.\nIf the print data is rotated 90 degrees to the vertical print orientation in the page mode, the desired print output cannot be achieved if one page of data in the printing area is distributed horizontally to the front and back sides. For example, when the data is converted for printing vertically, lines that are longer than the width of the coupon or receipt can be printed in the paper feed direction to one side of the print medium, but if the same data is printed horizontally, the printout will be split."}
{"text": "Various methods are currently used to obtain matt or satin powder paints but, to a greater or lesser extent, all of them have limitations.\nSome methods have recourse to overpigmentation or the addition of polymers which are incompatible with the binder system in intimate homogenization but this leads to a diminution in the chemical and mechanical resistance of the paints.\nOther methods use intimate homogenization of additives or special catalysts which, however, can only be used with specific systems of binders.\nOther methods, always with intimate homogenization, are based on the difference in reactivity between resins of the same type of reactive group in a single binder system.\nStill other methods do not use intimate homogenization but cold physical mixing of one or more paints extruded separately and with different viscosity or reactivity or systems of binders as described in international patent application PCT/EP 89/00022 in the name of the same applicant, published under number WO 89/06674.\nThese last methods are much less limiting than the preceding ones because they can employ practically all the systems of binders as long as they have different viscosities or different types of cross linking. However they are less practical because, with cold mixing methods, manufactures of powder paints normally extrude the paints separately, cool the extrusions, crush or grind them, sieve them and then weigh them out in the various pre-established ratios and mix them in a mixing machine. This means that the process is discontinuous and thus more time, space, materials and manpower are required.\nThere is also the danger of partial separation of the two paints during the various stages of grinding and sieving because of the difference in their specific weights or a different conformation of the granules of powder. Finally, at the spraying stage, the two paints can be deposited non-uniformly because of the different charges accumulated by the two binders."}
{"text": "Device shrinking requires scaling down the supply voltage (Vcc) accordingly because of power budget limitations. However, bit cells may require a longer time to reach a stable state after a write operation at lower supply voltages. This presents a trade-off between the minimum supply voltage (Vccmin) and the read/write speed of a memory.\nMemories (e.g., cache memories, instruction queues, translation look-aside tables) consume a significant amount of switching power and leakage power in chips. Some bit cell designs that can operate at lower supply voltages incur overhead in terms of delay."}
{"text": "More than one million micro-electronic field effect transistors (e.g., complementary metal-oxide-semiconductor (CMOS) field effect transistors) are fabricated on a substrate (e.g., semiconductor wafer) to form an integrated circuit. During the fabrication process of the integrated circuit, a photoresist is deposited, exposed, and developed to create a mask utilized to etch the underlying layers.\nTo produce the integrated circuit, implanting ions into various portions of the integrated circuit is necessary. During ion implantation, wafers are bombarded by a beam of electrically charged ions, called dopants. Implantation changes the properties of the material where the dopants are implanted in order to achieve a particular electrical performance. These dopants are accelerated to an energy that will permit them to penetrate (i.e., implant) the film to a desired depth. During implantation, high-dose or high-energy ions may implant in the photoresist layer and cause a hard, crust-like layer to form on the surface of the photoresist layer. The crust layer is difficult to remove using conventional stripping processes. Moreover, if the crust layer or underlying photoresist is not removed, the residual photoresist becomes a contaminant during subsequent processes. Thus, an improved method for stripping a photoresist is still in great demand."}
{"text": "In poker, baccarat, bridge, blackjack, and other card games, a dealer sets one or more decks of playing cards in a card shoe or the like, and deals cards to game players by drawing them therefrom one by one. In so doing, to ensure the fairness of the games, the cards need to be dealt at random. Therefore, a game host must sufficiently shuffle the playing cards randomly to ensure a random order of arrangement of the playing cards before they are set in the card shoe. The shuffled playing cards used in various card games such as poker, baccarat, bridge or blackjack will ordinarily include 416 cards if eight decks of cards are used.\nA conventional card shuffling device for shuffling cards is disclosed in, for example, Patent Literature 1. Each shuffled playing card set is provided after being shuffled with a shuffling apparatus to be arranged in a unique order with a uniquely identifiable shuffled card ID affixed on its packing box or the like as a barcode.\nAlso, a technique for printing on each card used in a casino or the like the suit (type) and the rank (or number) of the card as a code in order to prevent counterfeiting is disclosed in Patent Literature 2 or the like.\nShuffled playing cards are dealt from a card shoe apparatus to be used in a game, and the code of each card dealt that represents its rank is read by the card shoe apparatus. In recent years, technology for counterfeiting cards has greatly advanced, and accordingly, illegal counterfeiting of cards with such code attached is no longer impossible. In a case where card itself is not falsified, a card which has regular code (qualified code) is illegally obtained. In such a case, before using the card for the game, this card can be mixed into a predetermined position among the shuffled playing cards. Further, in such a case, the falsified shuffled playing cards which is intentionally ordered can be wholly replaced for regular cards (8 decks of cards)."}
{"text": "Organizations such as on-line retailers, Internet service providers, search providers, financial institutions, universities, and other computing-intensive organizations often conduct computer operations from large scale computing facilities. Such computing facilities house and accommodate a large amount of server, network, and computer equipment to process, store, and exchange data as needed to carry out an organization's operations. Typically, a computer room of a computing facility includes many server racks. Each server rack, in turn, includes many servers and associated computer equipment.\nComputer systems typically include a number of components that generate waste heat. Such components include printed circuit boards, mass storage devices, power supplies, and processors. For example, some computers with multiple processors may generate 250 watts of waste heat. Some known computer systems include a plurality of such larger, multiple-processor computers that are configured into rack-mounted components, and then are subsequently positioned within a rack system. Some known rack systems include 40 such rack-mounted components and such rack systems will therefore generate as much as 10 kilowatts of waste heat. Moreover, some known data centers include a plurality of such rack systems.\nComputer systems typically include a number of components that generate waste heat. Such components include processors, mass storage devices, and power supplies. For example, some computers with multiple processors may generate 250 watts of waste heat. Typically, each chassis in the rack has air vents in the front and rear panels. An external cooling system causes air to flow through the front vents, over heat producing component in the rack, and out through the back vents. In many cases, heat sinks, heat pipes or other components are required to achieve an adequate cooling of the heat producing components of a system. External cooling systems, thermal management components, and the associated parts such as mounting brackets, spring clips, and fasteners, add cost and complexity to a system.\nIn a typical computing device, the highest heat-generating components and most sensitive components are located at a few locations in the device. Thus, cooling needs are not evenly distributed over the computing device, but are instead concentrated at particular locations. The location of the highest heat-generating components in rack mounted computing devices (such as rack-mounted servers) varies across different models of computing devices. For example, in one rack-mounted server model, the central processing units may be in the front left corner of the server, while in another rack-mounted server model, the central processing unit may be in the right rear corner of the server.\nHard disk drives include motors and electronic components that generate heat. Some or all of this heat must be removed from the hard disk drives to maintain continuous operation of a server. The amount of heat generated by the hard disk drives within a data room may be substantial, especially if all of the hard disk drives are fully powered up at all times.\nWhile the invention is susceptible to various modifications and alternative forms, specific embodiments thereof are shown by way of example in the drawings and will herein be described in detail. It should be understood, however, that the drawings and detailed description thereto are not intended to limit the invention to the particular form disclosed, but on the contrary, the intention is to cover all modifications, equivalents and alternatives falling within the spirit and scope of the present invention as defined by the appended claims. The headings used herein are for organizational purposes only and are not meant to be used to limit the scope of the description or the claims. As used throughout this application, the word “may” is used in a permissive sense (i.e., meaning having the potential to), rather than the mandatory sense (i.e., meaning must). Similarly, the words “include,” “including,” and “includes” mean including, but not limited to."}
{"text": "1. Field of the Invention\nThis invention relates generally to rotary drill bits used in drilling subterranean wells and, more specifically, to rotary drill bits employing a torque limiting device allowing the drill string to rotate relative to the crown of the bit when a predetermined reactive torque is experienced by the crown of the drill bit.\n2. State of the Art\nThe equipment used in drilling operations is well known in the art and generally comprises a drill bit attached to a drill string, including drill pipe and drill collars. A rotary table or other device such as a top drive may be employed to rotate the drill string, resulting in a corresponding rotation of the drill bit. The drill collars, which are heavier and stiffer than drill pipe, are normally used on the bottom part of the drill string to add weight to the drill bit. The weight of these drill collars assists in stabilizing the drill bit against the formation at the bottom of the borehole, causing it to drill when rotated. Too much weight on bit (WOB), however, may cause the drill bit to stall.\nDownhole motors may also be employed to rotate the drill bit and include two basic components: a rotor, which is a steel shaft shaped in the form of a spiral or helix, and a stator, which is a molded rubber sleeve in a rigid tubular housing, that forms a spiral passageway to accommodate the rotor. When the rotor is fitted inside the stator, the difference in geometry between the two components creates a series of cavities through which drilling fluid is pumped. In doing so, the fluid displaces the rotor, forcing it to rotate as the fluid continues to flow between the rotor and the stator. An output shaft connected to the rotor transmits its rotation to the bit.\nA typical rotary drill bit includes a bit body secured to a steel shank having a threaded pin connection for attaching the bit body to the drill string or the output shaft of a downhole motor and a crown comprising that part of the bit fitted with cutting structures for cutting into an earth formation. Generally, if the bit is a fixed-cutter or so-called \"drag\" bit, the cutting structure includes a series of cutting elements made of a superabrasive substance, such as polycrystalline diamond, oriented on the bit face at an angle to the surface being cut. On the other hand, if the bit has rotating cutters such as on a tri-cone bit, each cone independently rotates relative to the body of the bit and includes a series of protruding teeth, which may be integral with the cone or comprise separately formed inserts.\nThe bit body of a drag bit is generally formed of steel or a matrix of hard particulate material such as tungsten carbide infiltrated with a binder, generally of copper-based alloy. In the case of steel body bits, the bit body is usually machined from round stock to the shape desired, usually with internal watercourses for delivery of drilling fluid to the bit face. Topographical features are then defined at precise locations on the bit face by machining, typically using a computer-controlled, five-axis machine tool. For a steel body bit, hardfacing may be applied to the bit face and to other critical areas of the bit exterior, and cutting elements are secured to the bit face, generally by inserting the proximal ends of studs on which the cutting elements are mounted into apertures bored in the bit face. The end of the bit body opposite the face is then threaded, made up and welded to the bit shank.\nIn the case of a matrix-type drag bit body, it is conventional to employ a preformed so-called bit \"blank\" of steel or other suitable material for internal reinforcement of the bit body matrix. The blank may be merely cylindrical and tubular, or may be fairly complex in configuration and include protrusions corresponding to blades, wings or other features on the bit face. Other preform elements comprised of sand, or in some instances tungsten carbide particles, in a flexible polymeric binder may also be employed to define internal watercourses and passages for delivery of drilling fluid to the bit face, as well as cutting element sockets, ridges, lands, nozzle displacements, junk slots and other external topographic features of the bit. The blank and other preforms are placed at appropriate locations in the mold used to cast the bit body before the mold is filled with tungsten carbide. The blank is bonded to and within the matrix upon cooling of the bit body after infiltration of the tungsten carbide with the binder in a furnace, and the other preforms are removed once the matrix has cooled. The threaded shank is then welded to the bit blank. The cutting elements (typically diamond, and most often a synthetic polycrystalline diamond compact, or PDC) may be bonded to the bit face by the solidified binder subsequent to furnacing of the bit body. Thermally stable PDCs, commonly termed \"TSPs\", may be bonded to the bit face by the furnacing process or may be subsequently bonded thereto, as by brazing, adhesive bonding, or mechanical affixation.\nIn order for the cutting elements to properly cut the formation during a drilling operation, considerable torque is required to generate the necessary rotational force between the cutting elements and the formation under a WOB substantial enough to ensure an adequate depth of cut. The resultant or reactive torque on the bit from formation contact is translated through the drill string and must be overcome by the means used to rotate the drill string, such as a rotary table, top drive, or downhole motor. In some instances, such as drilling through harder formations, the resultant torque may result in the winding up and sudden release of the drill string under torque, manifested as so-called \"slaps\" of the drill string at the rotary table. In other instances, torque may be sufficient to actually stop the bit from rotating. The rotary table may continue to rotate the drill string for some time, in effect \"twisting\" the drill string and placing the bit under very high torque loads before an operator realizes that the bit is no longer rotating. This problem is of particular concern with drag bits, due to direct engagement of the formation by the fixed PDC cutters, but also manifests itself with rock bits. If such a condition occurs and the rotary table continues to rotate, the drill string, the bit and/or components thereof may be damaged, or the drill string may even part under the torque load. If failure of the drill string occurs, the portion of the drill string above the break must be removed from the wellbore. A \"fishing\" assembly inserted into the wellbore is then normally employed in an attempt to retrieve the remainder of the drill string. If retrieval is impractical or unsuccessful, a new drilling assembly must be deflected, \"sidetracked,\" or steered around the \"fish.\" Any such scenario adds to the cost of production and results in down-time of the drilling operation while the remainder of the broken drill string is \"tripped\" from the wellbore and replaced with other bottom hole assemblies.\nWhen a downhole motor is being used to rotate the drill bit, a sudden rise in surface pressure of the drilling fluid may indicate that the motor has stalled. While other conditions may cause a rise in fluid pressure, such as a clogged motor or plugged nozzles, if the motor stalls because the bit is no longer rotating due to excessive torque on the bit and is maintained in a stalled condition, the elastomeric stator lining may be damaged, preventing a proper interface between the stator and the rotor, thus requiring the motor to be tripped out of the wellbore and replaced. At the least, the bottomhole assembly, including the motor, must be pulled off-bottom and drilling and circulation recommenced to start the motor before the formation is re-engaged by the bit.\nIn addition to damage to drill strings and bits, directional drilling presents its own set of problems when excessive torque is applied to the drill bit. A directional well must intersect a target that may be several miles below the surface location of the drilling rig, and laterally offset therefrom. In order to reach the target, the wellbore must be directed or steered along a predetermined trajectory. The trajectory of the bit is typically determined by the tool face orientation (TFO), which must be maintained during drilling in order to maintain the trajectory of the wellbore toward the desired target. If the TFO shifts due to a stalled drill bit, the drilling must stop and a new TFO set as a reference point for the direction of drilling. While a shift in TFO is quickly manifested to the operator due to the essentially real-time nature of the MWD (measurement while drilling) mud-pulse transmissions, nonetheless, loss of TFO and resetting thereof results in considerable reduction in the overall rate of penetration (ROP) of the drilling assembly.\nIt would thus be advantageous to provide a drill bit assembly that includes a torque limiting device that is either an integral part of the bit construction or is attached near the bit between the drill bit and the drill string, or is positioned between the downhole motor and the drill bit."}
{"text": "1. Field of the Invention\nExemplary embodiments of the present invention relate to a graphical user interface (GUI) for electronic devices and, in particular, to a method and apparatus for displaying a GUI according to a user's contact pattern.\n2. Description of the Background\nA touch screen may be used as a display unit and an input unit. Therefore, electronic devices having a touch screen may not require additional display and input units. Because of this advantage, a touch screen may be widely used for limited-size electronic devices such as, for example, mobile devices, which may also be referred to as portable devices or handheld devices.\nTypically, a user can manipulate a touch screen with one hand or both hands to command an execution of a desired function or application. When a user uses both hands, one hand may grip a device, and the other hand may touch a touch screen of the device. However, when a user uses only one hand, parts of the touch screen may often be hidden by a finger (e.g., the thumb) of a gripping hand.\nFIG. 10A is an illustrative example showing a user's left thumb selecting one of the menu icons displayed on a touch screen. In this example, if a user touches a specific icon (e.g., Music Icon) located in the upper right portion on the touch screen, some of the other icons (e.g., Game Icon, Display Icon, and Scheduler Icon) displayed on the touch screen may be entirely or partially hidden by the thumb. Additionally, these hidden icons may be in contact with the thumb, and thereby functions associated with the hidden icons may be undesirably executed.\nFIG. 10B is another illustrative example showing a user's left thumb touching a scroll bar presented on the touch screen. If the user touches the scroll bar located at the right side on the touch screen, displayed content (e.g., scene) may be hidden by the thumb. Furthermore, some displayed content may undesirably be touched with the thumb and accessed.\nAn electronic device having tactile sensors may provide control of electronic device applications only if a user keeps contact with a specific part of the device without using a touch screen or a keypad. Such electronic devices may provide a display screen with a GUI to guide contact-based inputs. If the GUI is displayed in a fixed form regardless of a user's contact pattern, some positions in the GUI may fail to note/input the user's contact. This may be due to differences in individual hand sizes, finger sizes, and grip forms. It is difficult to realize a GUI adapted for multiple users. If no match exists between a position in the GUI and a contact point of the user's contact, confusion may be generated when a user manipulates applications on the electronic device."}
{"text": "This specification relates generally to sharing learned information among robots.\nRobots can be configured to perform a variety of tasks. Robots generally need some form of perception to complete a task, such as the recognition or identification of an object. Robots may encounter objects and situations that their programming or training does not prepare them to handle appropriately. For example, although a robot may be able to recognize many different types of objects, it is likely that the robot will likely encounter new objects that it does not recognize."}
{"text": "1. Field of the Invention\nThe invention relates to a motor element, in particular to a dental medical handpiece.\n2. Description of the Related Technology\nA handpiece of this kind is described in DE 44 06 855 A1. In the case of this previously known handpiece, a tool-holder that is formed as a collet chuck is pushed into a drive shaft that is formed so that it is hollowly cylindrical at least on the front side, with a coupling for the axial connection of the tool-holder to the drive shaft being formed by a screw coupling. Consequently, the tool-holder is connected to the drive shaft by screwing in and can be disconnected again by unscrewing, in order to clean it for example. This outlay on assembly and disassembly is comparatively great. Moreover, it is mechanically difficult to fix the tool-holder in the screwed-in state in order to avoid unintentional untwisting. It is to be taken into consideration in this connection that during the functional operation of the handpiece when switching the drive shaft on and off torques result that can loosen a screw securing arrangement based on screw-fastening so that the function of the collet chuck can be impeded and impaired without this defect being directly noticeable. A further disadvantage lies in the fact that screw-fastening is comparatively difficult, because the collet chuck is arranged so that it is largely sunk in the drive shaft and it is difficult therefore to reduce a torque for screw-fastening at the front end of the collet chuck.\nThe underlying object of the invention is therefore to improve the coupling of the tool-holder or the locking and/or unlocking of the coupling in the case of a handpiece in accordance with the invention."}
{"text": "1. Field\nAspects of the present invention relate to a rechargeable battery. More particularly, aspects of the present invention relate to a rechargeable battery having a short circuit member.\n2. Description of the Related Art\nUnlike a primary battery, a rechargeable battery is designed to repeatedly perform charge and discharge functions. Small capacity rechargeable batteries may be used in small portable electronic devices such as mobile phones, laptop computers, and camcorders, and large capacity rechargeable batteries may be used in power sources for driving motors of hybrid vehicles and the like.\nIn recent years, a high power rechargeable battery using a non-aqueous electrolyte having high energy density has been developed. The high power rechargeable battery is a large capacity rechargeable battery with a plurality of rechargeable batteries connected to each other in series to be used to drive a motor of a device requiring high power, such as an electric vehicle.\nThe large capacity rechargeable battery generally includes a plurality of rechargeable batteries which are connected to each other in series, and the rechargeable battery may have a cylindrical shape or an angular shape.\nSuch a rechargeable battery requires a safety apparatus capable of reducing or preventing excessive heat and ensuring safety upon overcharge.\nThe above information disclosed in this Background section is only for enhancing understanding of the background of the invention and therefore it may contain information that does not form prior art."}
{"text": "Fuel cell systems are currently being developed for use as power supplies in numerous applications, such as automobiles and stationary power plants. Such systems offer promise of economically delivering power with environmental and other benefits.\nFuel cells convert reactants, namely fuel and oxidant, to generate electric power and reaction products. Fuel cells generally employ an electrolyte disposed between two electrodes, namely a cathode and an anode. A catalyst typically induces the desired electrochemical reactions at the electrodes. Preferred fuel cell types include solid polymer electrolyte fuel cells that comprise a solid polymer electrolyte and operate at relatively low temperatures.\nDuring normal operation of a solid polymer electrolyte fuel cell, fuel is electrochemically oxidized at the anode catalyst, typically resulting in the generation of protons, electrons, and possibly other species depending on the fuel employed. The protons are conducted from the reaction sites at which they are generated, through the electrolyte, to electrochemically react with the oxidant at the cathode catalyst. The catalysts are preferably located at the interfaces between each electrode and the adjacent electrolyte.\nA broad range of fluid reactants can be employed in solid polymer electrolyte fuel cells and may be supplied in either gaseous or liquid form. For example, the oxidant stream may be substantially pure oxygen gas or a dilute oxygen stream such as air. The fuel may be substantially pure hydrogen gas, a gaseous hydrogen-containing reformate stream, or an aqueous liquid methanol mixture in a direct methanol fuel cell. Reactants are directed to the fuel cell electrodes and are distributed to catalyst therein by means of fluid diffusion layers.\nSolid polymer electrolyte fuel cells employ a membrane electrode assembly (“MEA”), which comprises the solid polymer electrolyte or ion-exchange membrane disposed between the two electrodes. Each electrode contains a catalyst layer, comprising an appropriate catalyst, located next to the solid polymer electrolyte. The catalyst may be, for example, a metal black, an alloy or a supported metal catalyst, for example, platinum on carbon. The catalyst layer typically contains ionomer, which may be similar to that employed for the solid polymer electrolyte (for example, Nafion®). The catalyst layer may also contain a binder, such as polytetrafluoroethylene. The electrodes may also contain a substrate (typically a porous electrically conductive sheet material) that may be employed for purposes of mechanical support and/or reactant distribution, thus serving as a fluid diffusion layer.\nThe MEA is typically disposed between two plates to form a fuel cell assembly. The plates act as current collectors and provide support for the adjacent electrodes. The assembly is typically compressed (for example, of order of 70 pounds per square inch (psi) overall) to ensure good electrical contact between the plates and the electrodes, in addition to good sealing between fuel cell components. A plurality of fuel cell assemblies may be combined in series or in parallel to form a fuel cell stack. In a fuel cell stack, a plate is usually shared between two adjacent MEAs, and thus also serves as a separator to fluidly isolate the fluid streams of the two adjacent MEAs.\nFurther, flow fields are typically incorporated into both surfaces of such plates in order to direct reactants across the electrochemically active surfaces of the fluid diffusion electrodes or electrode substrates. The flow fields typically comprise fluid distribution channels separated by landings. The channels provide passages for the distribution of reactant to the electrode surfaces and also for the removal of reaction products and depleted reactant streams. The landings act as mechanical supports for the fluid diffusion layers in the MEA and provide electrical contact thereto. Ports and other fluid distribution features are typically formed in the surfaces at the periphery of the flow field plates. When assembled into a fuel cell stack, the stacked ports can form internal manifolds for distribution of the fluids throughout the stack. The other distribution features typically are provided to distribute fluids from the ports to the appropriate flow fields. PCT/International Publication No. WO 00/41260, for instance, illustrates flow field plates with flow fields comprising a plurality of straight, parallel channels. The flow fields are fluidly connected to manifold openings in header regions at the periphery of the plates by a series of complex passages formed in the plate surfaces.\nSealing of some sort is generally required around the edges of the MEAs to isolate the different fluids on each side of the MEA. In principle, the membrane electrolyte in the MEA may be oversized and extend significantly beyond the electrochemically active area in order to be employed as a sealing gasket. However, membrane electrolyte material is generally expensive and has relatively poor mechanical properties for this purpose. Thus, it is preferred to employ other means for effecting edge seals.\nAs an alternative, U.S. Pat. No. 5,464,700 discloses a gasketed membrane electrode assembly that employs gasketing material at the membrane periphery, rather than the membrane itself, as a gasket. The gasketing material may be formed from an elastomeric material suitable for cold bonding or bonding by heat and pressure. A non-hydrophilic thermoplastic elastomer is the preferred gasketing material (for example, Santoprene brand gasketing material).\nFurther, U.S. Pat. No. 5,187,025 discloses a unitized cell assembly in which the edge of the MEA is extended with a laminated plastic structure to give it strength and rigidity for sealing and support. Therein, the electrolyte membrane is surrounded with a plastic spacer having a thickness closely matched to that of the membrane. A thin, plastic film is bonded with an adhesive layer to both sides of the membrane and spacer so that the film and adhesive bridge the gap therebetween. Porous electrodes with plastic frames are bonded to the composite membrane. The use of this type of structure permits the construction of a rigid cell frame, which can be made the same thickness as the membrane electrode package.\nIn an effort to further improve fuel cell performance and to reduce the thickness and cost of fuel cell assemblies, there is a trend to use thinner components and especially thinner flow field plates. A particularly thin design for flow field plates employs corrugated flow fields which are characterized by features on one side that are complementary with features on the other side. Corrugated structures are readily formed out of thin metallic sheets by stamping or rolling methods. However, it is not so easy to use such methods to form features on opposing sides of such sheets if the intended features are not substantially complementary. (Instead, engraving techniques may be employed on thicker starting sheet material.) Yet, some desirable corrugated flow field designs involve complex flow distribution and mechanical support features at the fluid inlet and outlet ports on one or both sides of the plate which, in some cases, cross over similarly complex features on the opposite side of the plate. Further, some designs may involve features that interconnect channels within the flow field itself (for example, which interconnect parallel flow channels at their ends so as to form serpentine channels). It might not be possible to form such features on both sides of a plate in a complementary fashion, thereby complicating the fabrication of plates with corrugated flow fields. For this and other reasons, it would be advantageous to be able to design and fabricate the components that distribute fluids to and from the flow fields more independently of the components that comprise the flow fields themselves."}
{"text": "The present application relates to the field of metrology, and particularly to LiDAR (Light Distance and Ranging) sensors and related devices.\nThe process for measuring distance and reflectance values for objects within an environment without touching those objects is of great interest to many industries including surveying, architecture, entertainment (character generated effects for movies and video games), construction, forensic and geography applications. Historically to collect accurate distance and reflectance measurements one used photogrammetry techniques, but the process for extracting information from stereo imagery is both time consuming and expensive. Over the past decade, advances in Light Detecting and Ranging (LiDAR) technology have enabled practitioners to scan large area surfaces while collecting billions of data points, each with a precise latitude, longitude and elevation (x, y, z) values within the local (relative) coordinate system. The aggregate of the billions of data points is referred to as a point cloud data set. Practitioners will subsequently extract objects from the point cloud and then create three dimensional models. Those models are then used in numerous applications. For example, within geographic information systems (GIS) industry, practitioners will frequently integrate Global Positioning System (GPS) data into the point cloud to ‘geo-reference’ it to a global coordinate system. Every data point in a geo-referenced point cloud has an absolute x, y, and z location on the earth's surface."}
{"text": "In the manufacture of integrated circuits from semiconductor wafers and substrates, wafer polishing is a common technique utilized to remove material and/or achieve planarity. Such polishing can be conducted by purely chemical, purely mechanical or chemical-mechanical polishing means (CMP). With CMP, polishing and removal occurs by a combination of both chemical and mechanical polishing action. CMP utilizes a combination of solid abrasives and chemicals to achieve the combination polishing action. One type of chemical-mechanical polishing utilizes a slurry comprising very hard, solid abrasive particles suspended in a chemical solution. The slurry is interposed between a pad and a wafer, with both typically being caused to rotate, and material removed from the wafer by both chemical and mechanical action. Another form of CMP provides abrasive material embedded within the surface of the polishing pad, and is commonly referred to as fixed abrasive CMP.\nUnfortunately, conventional CMP slurries designed for non-fixed abrasive CMP create problems and do not always work satisfactorily in fixed abrasive CMP processes. This has been discovered to be particularly true in the CMP of layers where copper is present at greater than or equal to fifty atomic percent. Increasing interest is being focused on copper as a next generation material for interconnect lines in semiconductor circuitry fabrication. CMP of copper in processes of forming such lines will play a significant and valuable role in the fabrication of such circuitry. While CMP of copper has been reported using non-fixed abrasive pads and slurries, existing materials have proven less than satisfactory when using fixed abrasive pads. Accordingly, needs remain for improved chemical-mechanical processes using fixed abrasive pads and in the development of polishing solutions therefor."}
{"text": "This invention relates in general to speech diagnosis therapy and, more specifically, to a tongue force measuring device for use in speech therapy.\nSome speech pathologists feel that certain speech disorders stem from poor tongue strength in one or more directions of tongue movement. Treatment of the problem is usually attempted when the subject is a child. Attempts have been made to measure tongue strength using a load cell placed toward the middle of the tongue with the head restrained. Since the cell measured only mid-tongue strength in a single direction, the results were of limited value. This work was described by J. P. Working et al. in their article \"Tongue Strength Part I: Following Total Laryngectomy\" The Laryngoscope, 90 (1980), pp. 680-684. Other attempts to measure tongue strength use head restraints and cumbersome equipment which tend to frighten small children, detracting from validity of results.\nA number of devices have been developed to detect linguapalatal contact during phonation by means of a pattern of spaced electrical contacts on an artificial palate. Typical of these are the devices described by Takinishi et al. in U.S. Pat. No. 4,310,002, Fletcher et al. in U.S. Pat. No. 4,112,596 and Hori in U.S. Pat. No. 4,287,895. While these provide information as to tongue location during speech which is useful in treating speech disorders, they provide no information on tongue strength or force in any particular direction during speech.\nAn oral strain gage is described by Davis et al. in U.S. Pat. No. 3,297,021. However, this device measures only pressure between the jaws during chewing or the like. Since the strain gage is embedded in an artificial tooth, it is incapable of measuring tongue force.\nThus, there is a continuing need for a device for measuring tongue force during speech in three dimensions while simultaneously being light, compact and suitable for use by children."}
{"text": "Information retrieval systems, such as databases or search engines, allow users to retrieve information related to a specific subject by using one or more keywords that may be related to the specific subject. For example, a legal search service called Lexis® provided by the LexisNexis Group is widely used in the legal field to search for cases, journal articles, treaties, as well as other publications that are related to a specific topic or issue. Another information retrieval system, Google®, provided by Google, Inc., is a search engine commonly employed by internet users to search for web sites or online documents that are related to a specific subject matter.\nIn order to search for and retrieve documents related to a specific subject matter, users need to formulate query that typically comprises a set of keywords, phrases, symbols, commands, and/or other entities that are considered to be relevant to the subject matter or possibly contained in the documents relating to the subject matter. This type of information retrieval system poses problems to users because the users need to be familiar with the proper format for inputting queries into such systems. In addition, users need to have a basic understanding of the subject matter to be searched as well as of properties of the language used to describe that subject in order to formulate proper query to conduct the search.\nSome information retrieval systems provide assistance on query formation. For example, a website www.ask.com provides a search function called Ask Jeeves that allows users to input their questions in natural language. The system will extract keywords from the questions and conduct a search accordingly. Lexis® also provides a similar function allowing users to input search terms in natural language, either as a question or a statement. The system then extracts keywords from such natural language inputs to search for information related to the keywords.\nAlthough these tools provide basic assistance on query formation during information search and retrieval, such tools cannot function effectively in more realistic work environments in which the content of the query of question plays a paramount role. For example, consider the following search scenarios related to the same keyword “caterpillar:”\nScenario 1:                A biology student writing a term paper on animal development. In this case, the information search should be related to metamorphosis, the process by which the caterpillar becomes a butterfly.        \nScenario 2:                A contractor working on a construction plan for a new building. The contractor is most likely referring to Caterpillar, Inc., a major manufacturer of construction equipment.        \nScenario 3:                A grade-school student writing a book report on Lewis Carroll's book, Alice's Adventures in Wonderland. In this case, information retrieved should preferably be related to the character in the book, chapter excerpts, and pictures that the student could include in her paper.        \nThese scenarios illustrate various problems associated with conventional information retrieval systems. The first problem is that conventional information retrieval systems do not consider relevance of active goals in searching for information. The active goals of the user contribute significantly to the interpretation of the search terms and to the criteria for judging a resource as being relevant to the search terms. Typically, these goals are not fully expressed by users in forming their queries when using conventional information retrieval systems.\nThe second problem is that conventional information retrieval systems are subject to word-sense ambiguity. For example, The word “caterpillar” in scenario 1 should be treated differently from that in scenario 2. The context of the request provides a clear choice of word sense between the insect and the company. Conventional information retrieval systems cannot distinguish the subtle differences unless additional keywords or information are provided by the user.\nThe third problem is that conventional information retrieval systems fail to consider audience appropriateness when searching and retrieving information based on keywords or queries provided by the user. In addition to the keywords provided by the user, attributes related to the user in each of above the scenarios should also influence the choice of results. Sources appropriate for an advanced biology student will likely not be appropriate for a student in grade school.\nMoreover, when using conventional information retrieval systems, users often are unable to provide sufficient information in their queries. Studies show that on average, users' queries tend to be two to three words long. Needless to say, a two-word query most likely does not contain enough information to discern the active goals of the user, or even the appropriate senses of the words in the query.\nFurthermore, even if the user has sufficient knowledge to formulate workable queries to conduct a search, the user must be aware of the variety of available resources, decide where to find them, and must know how to use different information retrieval systems correctly, including details such as those concerning special operators like “and,” “or,” or “+” that are used differently in different information retrieval systems.\nTherefore there is a need to provide an automatic query formation system to assist users in retrieving information related to their active goals without their intervention. There is another need for an information retrieval system to consider the context of words or phrases when conducting an information search and retrieval. There is also a need to improve the performance of an information retrieval system by refining queries based on various attributes related to the users. An additional need exists to automate the information search and retrieval process by forming queries in proper format for conducting information search in different information sources."}
{"text": "1. Field of the Invention\nThe present invention relates to an inputting method and an input device suitable to perform an input operation surely by pressing of an operator.\n2. Description of the Related Art\nA flat input device is used as an input device for electric devices such as a mobile phone and a personal digital assistant (PDA).\nAs an example of such an input device, there is given a flat input device including a capacitance sensor or a pressure sensor and capable of coordinate input. In the input device including a capacitance sensor or a pressure sensor (hereinafter referred to simply as “sensor”), input of a numerical value or a symbol is achieved with just a single touch of the sensor. Therefore, an operator has no sensation of performing an input operation and unintended inputting errors tend to occur.\nTo solve this problem, an input device is disclosed. The input device has a switch in addition to a sensor. The switch is disposed behind the sensor. The switch responds to a pressing operation of an operator by reaction force. The input device thus makes the operator realize that an input operation is performed, thereby preventing inputting errors. The input device is disclosed in, for example, Japanese Unexamined Patent Application Publication No. 2002-123363.\nAlthough the known input device can make an operator realize that an input operation is performed, it has a problem in which what is input is not known until the operator looks at a display displaying the input result (confirmed input) as data (output of the input device).\nThat is to say, the known input device has a problem in which it is impossible to know whether the input operation is correct or not before looking at the displayed data and therefore an inputting error can occur. The inputting error is that the region actually pressed is different from the intended region, for example, that an adjacent region is pressed by mistake. The smaller the operating portion is, the more often such an inputting error occurs. In the case where such an inputting error occurs, it is necessary to perform a corrective input operation, that is to say, to delete the input data and then to input it again. This is less user-friendly.\nTherefore, an inputting method and an input device with which an operator can perform a correct input operation surely and easily and which are user-friendly are needed."}
{"text": "A common procedure for handling pain associated with intervertebral discs that have become degenerated due to various factors such as trauma or aging is the use of intervertebral fusion devices for fusing one or more adjacent vertebral bodies. Generally, to fuse the adjacent vertebral bodies, the intervertebral disc is first partially or fully removed. An intervertebral fusion device is then typically inserted between neighboring vertebrae to maintain normal disc spacing and restore spinal stability, thereby facilitating an intervertebral fusion.\nThere are a number of known conventional fusion devices and methodologies in the art for accomplishing the intervertebral fusion. These include screw and rod arrangements, solid bone implants, and fusion devices which include a cage or other implant mechanism which, typically, is packed with bone and/or bone growth inducing substances. These devices are implanted between adjacent vertebral bodies in order to fuse the vertebral bodies together, alleviating the associated pain.\nHowever, there are drawbacks associated with the known conventional fusion devices and methodologies. For example, present methods for installing a conventional fusion device often require that the adjacent vertebral bodies be distracted to restore a diseased disc space to its normal or healthy height prior to implantation of the fusion device. In order to maintain this height once the fusion device is inserted, the fusion device is usually dimensioned larger in height than the initial distraction height. This difference in height can make it difficult for a surgeon to install the fusion device in the distracted intervertebral space.\nFurther, lordosis refers to a curvature of the spine, and in particular a curvature that is posteriorly concave. In certain patients, this curvature may, for example, be larger than desired. Traditional vertebral fusion procedures and devices do not adequately account for this curvature. As such, traditional devices do not properly align with adjacent vertebral bodies. To ensure proper fit of traditional devices, bone may be removed from the vertebral bodies, increasing procedure and healing time.\nAs such, there exists a need for a fusion device capable of being installed inside an intervertebral disc space at a minimum distraction height and for a fusion device that can maintain a normal distance between adjacent vertebral bodies when implanted."}
{"text": "1. Field of the Invention\nThe invention relates to a lever-type connector, to a connector assembly and to an assembling method therefor.\n2. Description of the Related Art\nJapanese Unexamined Patent Publication No. 2004-14142 discloses a lever-type connector with first and second housings that are connectable with each other. A lever formed with a cam groove is mounted in first housing and the second housing includes a cam pin. A jaw bulges radially out at the leading end of the cam pin and engages a step at an intermediate position of the cam groove. The step extends substantially horizontally at the same depth along the cam groove, and a space above the step accommodates the jaw. The cam pin enters the cam groove when the two housings are fit lightly together. The lever then is rotated so that the cam pin moves towards the back of the cam groove to urge the housings to a properly connected state. The jaw slides on the step so that the lever does not come out of the first housing during the connecting operation.\nA clearance is formed between the step of the cam groove and the jaw due to a dimensional tolerance or the like. If this clearance is too small, it may be difficult for the jaw to move above the step. Thus, a contact pressure of the step of the cam groove and the jaw becomes too high, which might result in difficulty in starting the lever rotation. On the other hand, if this clearance is too large, the central axis of the cam pin may incline with respect to a vertical axis upon the receipt of connection resistance of the two housings while the cam pin is moving in the cam groove. In short, there is no problem if the clearance is set suitably, but it is difficult to constantly define a specified clearance for each product in view of variation in production.\nThe invention was developed in view of the above situation, and an object thereof is to improve the operability of a movable member for connecting or assisting the connection of a connector with a mating connector."}
{"text": "1. Field of the Invention\nThis invention concerns the counting of items and simultaneous marking of said items during counting, and further relates to apparatus for accurately achieving said simultaneous counting and marking.\n2. Description of the Prior Art\nIn the course of conducting business activities, it is often necessary to take an inventory of items in storage as on shelves or in bins, racks or other storage systems. The counting of items of similar or identical appearance can be a tedious chore and subject to considerable error, particularly if the items are subject to movement in the course of the counting operation.\nTo enhance the accuracy of counting, one technique involves the marking of an item at the instant that it is being counted. Devices for the simultaneous marking and counting of items have earlier been disclosed, as for example in the following U.S. Patents:\nU.S. Pat. No. 4,993,050 to Carpenteri, et. al., discloses a counter/marker apparatus comprising a head assembly counter mechanism and attached first elongated sleeve. A second sleeve, portioned within said first sleeve for axially slidable reciprocal motion therein, is adapted to carry a marking pen. Relative movement between said sleeves causes activation of a switch which registers one unit of count upon said counter mechanism. A restoring spring urges the tip of the marking pen forwardly.\nAlthough the Carpenteri, et. al., apparatus is based upon sound general concepts, one of its shortcomings is that, as indicated at column 2, lines 50-53, the tip of the marking pen must be displaced rearwardly by at least 1/8 inch to activate the switch. The problem is that, if the pen is not displaced rearwardly by at least 1/8 inch, a mark will be made on an item, but no count will be made. It is also possible that, if the travel of the pen is not a pure reciprocal harmonic motion, as a result of a pause or axially transverse wiggle, more than one count may be registered with a single marking.\nU.S. Pat. No. 4,726,044 to Perna, et. al., discloses a counter/marker device in the form of a ball point pen having a movable ball point that is passed upwardly in the direction of the housing of the pen when a mark is made. Such upward movement closes a switch which activates an electronic counter. A second switch, activated by gravity, resets the counter for counting the next mark. The Perna, et.al., device will only function in marking substrates that are downwardly disposed. Also, there is no assurance that the counter is activated even though a mark has been made.\nU.S. Pat. No. 4,048,478 to Miwa, et. al., concerns a pen-shaped marking apparatus with electronic counter. The apparatus utilizes a single push button type of switch to activate the counter.\nU.S. Pat. No. 4,295,038 to Kreinbrink, et. al., discloses a pen-shaped marker connected by an electrical conductor to an otherwise separate counter device. As in the foregoing devices, there is no fail-safe provision to assure that, for each mark, there is a single count.\nU.S. Pat. No. 4,532,642 to Morris, et.al., discloses a counting stamper apparatus. A switch, mounted in the handle of the apparatus, is contacted by the upwardly directed stamping element to activate an electronic counting circuit.\nIt is accordingly an object of the present invention to provide a marking and counting device which will infallibly register a single count for each mark.\nIt is a further object of this invention to provide a device as in the foregoing object which can be employed in any attitude of use.\nIt is another object of the present invention to provide a device of the aforesaid nature which is of compact construction and easy to use.\nThese objects and other objects and advantages of the invention will be apparent from the following description."}
{"text": "1. Field of the Invention\nThe present invention is related to a device for transmission of force from the pistons of a piston engine, according to the preamble of the claims.\n2. Description of the Prior Art\nTransmission of force from the pistons of a piston combustion engine normally is made by the piston rods of the cylinders to a crankshaft. The transmission to the crankshaft is dependent on the angle of the crank arm in relation to the movement of the piston, thereby transforming the linear movement to the rotation of the crankshaft."}
{"text": "Many computing architectures implement cache as a means for addressing a mismatch between the performance characteristics of a computing device's storage device (e.g., magnetic recording media, solid state storage) and the performance requirements of the device's central processing unit (CPU) and application processing. For example, when an application issues a write command, the system may write the data into the cache, where periodically, the data from the cache is flushed to the storage device.\nHowever, flushing a write cache is an expensive operation, which can degrade performance. When flushing data from the cache to the storage device, the storage device flushes the entire system. That is, the storage device writes everything in the cache for all drives to storage, which can cause a computing device to incur a performance penalty, particularly on a low-end device with a limited bandwidth input/output (I/O) subsystem. Additionally, there may be a risk of loss or corruption of data if power is lost while there is data in the disk's write cache that has not been written to storage."}
{"text": "The present invention relates generally to damping devices for damping the vibration of overhead conductors, and particularly to a low-cost, corona-free damping device that takes the place of the well-known Stockbridge damper.\nThe inventor of the present disclosure is named in two previous patent applications directed to damping devices for overhead conductors that employ elastomer damping elements in a compact manner; namely, U.S. Ser. No. 565,180, filed Dec. 23, 1983 (now U.S. Pat. No. 4,527,008), and a continuation-in-part application thereof, U.S. Ser. No. 718,549, filed Apr. 1, 1985 (now U.S. Pat. No. 4,554,402). In addition, this application is being simultaneously filed with a second application by the inventor, the second application being entitled \"Damper With Stamped and Sand-Cast Parts,\" U.S. Ser. No. 805,819 filed Dec. 6, 1985.\nIt is well known that manufacturing companies look for economies in manufacture and improvements in performance of their products to meet competition. In the case of damping devices employing inertial weights attached to overhead conductors, the weights and clamp arms thereof are usually die cast. The machines and die casting molds employed to die cast components are costly, and in the case of overhead dampers, which are required to dampen the vibration of a multitude of conductor sizes (diameters), the cost is substantial. For example, there are about 350 conductor sizes presently employed in the United States to transmit electrical energy. All of these conductors are subject to aeolian vibration, and hence all require damping devices in one form or another. The damping devices require 21 clamp sizes to clamp the devices to the conductors. It can be appreciated that if separate die tooling is required for each of the conductor sizes, the cost can be prohibitive.\nTo properly match the impedance of the existing commercial conductors to damping devices, six inertial weight sizes are required. In the case where each damper required two, opposed weight halves, separate die tooling for each half would double the above cost of the weights."}
{"text": "In semiconductor device manufacturing, photolithography is typically used to transfer a pattern for forming semiconductor features onto the semiconductor wafer for the formation of multi-layered structures forming integrated circuits. During a photolithographic process, radiant energy having relatively small wavelengths such as ultraviolet light is passed through a photomask also referred to as a reticle to expose a radiant energy sensitive material such as photoresist formed on the wafer process surface. The mask includes predetermined circuitry patterns having attenuating regions and non-attenuating regions where the radiant energy intensity is modulated. For example, Ultraviolet (UV) light passed through the photomask onto the photoresist causes chemical reactions in the exposed portion of the photoresist altering it properties. Upon development of the photoresist resist exposed portions are removed in the case of a positive photoresist and non-exposed portions are removed in the case of a negative photoresist forming a pattern for subsequent processes such as anisotropic etching.\nAs semiconductor device feature sizes have decreased to sizes smaller than the wavelength of light used in photolithographic processes, the stray light incident on the exposed photoresist has increasingly become a problem in forming features with small critical dimensions (CDs), for example less than about 0.25 microns. Scattered light from undesired sources can cause a loss of pattern resolution in transferring the reticle pattern to the wafer photoresist. To increase the resolution of a transferred photolithographic pattern, phase shift masks (PSMs) have been developed where the phase of the wavefronts of light passing through alternating portions of the reticle pattern are shifted out of phase with respect to light passing through adjacent portions to produce destructively interfering wavefronts thereby reducing undesired exposure of the wafer photoresist due to diffraction of light at feature edges of the recticle pattern. As a result, the contrast, and therefore transferable resolution of the reticle pattern is improved.\nThere have been several different types of masks developed to improve resolution for different types of reticle patterns. For example, in an attenuated or halftone phase shift mask, the phase shifting function is typically accomplished by adding an extra layer of transmissive material to the mask with predetermined optical properties. Some PSMs are designed to produce improved resolution while having little improvement in depth of focus, while other PSMs are designed to have relatively modest increases in resolution while producing a greater improvement in depth of focus. For example, attenuated PSMs, also referred to as halftone PSMs, are of the latter type.\nIn a conventional mask, an opaque layer is formed which is not transmissive to light. Metals such as chromium are frequently used to form the opaque layer. The opaque layer is typically photolithographically patterned and etched, for example using an E-beam, laser or conventional UV light source to pattern a resist layer followed by etching to form a circuitry pattern including lines, pads, and contact holes.\nFor example, in the exposure process a mask formed on what is referred to as a reticle is repeatedly used to expose the circuit pattern on the reticle onto the photoresist covered process wafer. Typically, the term mask and reticle are now interchangeably used, although a photomask in the past has typically been used to refer to a mask that contains the pattern image for a complete wafer die array. Masks used in modern technology are typically referred to as reticles as well as masks where the reticle includes one or more mask patterns for individually patterning wafer die. In modern photolithography practice, a step and repeat process, such as step and scan, is used to expose the wafer surface with light passed through the mask over multiple wafer die. The mask on the reticle is larger by a factor of about 4 to 5 and is reduced to the appropriate dimensions on the wafer surface by optical reduction methods. For example, in the exposure process, the mask is illuminated by a light source either centered on the centerline of the projection optics or at an angle to the centerline of the projection optics, referred to as off-axis illumination to reduce the resolution limit and increase the depth of focus.\nOne problem with prior art processes is the presence of scattered light present during the exposure process. Scattered light present in the exposure process can alter a carefully designed exposure of the photoresist. For example, the light dose in an exposure process is critical for appropriate exposure of the photoresist to achieve desired pattern feature resolution, also referred to as critical dimension (CD), in the photoresist development process. Scattered light present in the exposure process can reduce the effectiveness of phase shift masks designed to reduce the effect of light diffraction around feature edges in the mask.\nPrior art processes have attempted reduce scattered light by the use of pellicle film, a thin layer of flexible and optically transparent material including anti-reflective properties that is tightly stretched on a frame about 5 mm to about 10 mm above the surface of the reticle. This configuration is not useful in many exposure processes, for example DUV exposure processes where the pellicle film may not be sufficiently transparent or may be degraded after repeated exposure.\nAnother shortcoming in the prior art relates to the thickness of the photoresist layer typically required to successfully etch PSMs, due to less than desirable selectivity in the etching process to the underlying PSM. As a result, as device sizes decrease it is difficult to achieve the desired CD even in the absence of undesired light reflections, as well as being affected in the exposure process by considerations of insufficient depth of focus.\nThus, there is a need in the semiconductor manufacturing art for an improved method and mask to reduce light scattering in a photolithographic exposure process while increasing an etching selectivity in a PSM etching process to improve critical dimension uniformity.\nIt is therefore among the objects of the present invention to provide an improved method and mask to reduce light scattering in a photolithographic exposure process while increasing an etching selectivity in a PSM etching process to improve critical dimension uniformity, in addition to overcoming other shortcomings and deficiencies of the prior art."}
{"text": "1. Field of the Invention\nThe invention relates to a front passenger airbag device.\n2. Description of the Related Art\nAn airbag door having four door panels that are separated from one another through the rupture of prearranged rupture portions extended in the shape of a combination of two letters “Y”, and provided with a regulation portion that is contiguous to a corner of an airbag passage port to regulate the direction of expansion of an airbag is described in Japanese Patent Application Publication No. 2006-62420 (JP-A-2006-62420).\nThe seating position of a passenger in a front passenger seat changes depending on the build of the passenger and the longitudinal position of the front passenger seat. It could be that sometimes the seating position of the passenger is sufficiently distant from a dashboard in which an airbag is accommodated or sometimes the seating position of the passenger is near a dashboard in which an airbag is accommodated.\nHowever, in the related art, the seating position of the passenger is not considered at all. In particular, there is still room for improvement in reduction of the reactive force of the airbag for the passenger when the passenger sits near the dashboard when being restrained by the expanded/deployed airbag."}
{"text": "In the realm of exercise equipment, punching and/or kicking bags and targets are commonplace. Existing stands are designed to withstand punching and kicks and as a result tend to be bulky and/or permanent fixtures. As such, existing stands are not easy to move to different locations and are not easy to store away when not in use.\nAdditionally, users tend to be of different heights and have different exercise needs (e.g., punching versus kicking exercises, martial arts versus boxing, etc.). However, many existing punching and kicking exercise equipment are not easily adjustable to accommodate these needs.\nOne type of prior art portable exercise device for training self defense includes a large hollow base that is filled with water or sand to make it stable when kicking or punching the pads or targets. Such devices are typically very heavy after filled with water and sand and as a result they are difficult to store or move from one location to another.\nAnother type of prior art portable exercise device for training self-defense includes a human torso and head shape for practicing punching or kicking techniques. Such devices are not adjustable in the height or position of the head and torso, and as a result may not be appropriately sized for very tall or short people, or practicing techniques for a variety of human sizes.\nConsequently, a portable stand for training self defense is needed that is stable, easily stored, moved around and easily adjustable over a range of heights and positions."}
{"text": "A battery state detection apparatus (battery sensor) that detects the state of a battery has been conventionally known. Such a battery state detection apparatus is disclosed in, for example, Patent Documents 1 and 2.\nThis type of battery state detection apparatus includes a circuit board and a shunt resistor for detecting a current. Referring to FIG. 9, a shunt resistor 7 is configured such that a resistor element 10 (made of Manganin, for example) whose resistance value is known is arranged between a first conductor part 11 and a second conductor part 12.\nEach of the first conductor part 11 and the second conductor part 12 is provided with a board connection terminal 15. The board connection terminal 15 includes a connecting part 30 to be connected to a circuit board (not shown).\nIn the conventional battery state detection apparatus, as shown in FIG. 9, the board connection terminal 15 is mounted to the shunt resistor 7 with a mounting screw. This structure, which enables the board connection terminal 15 to be electrically and mechanically connected to the shunt resistor 7 with reliability, can ensure a high degree of accuracy of sensing.\nThe circuit board is configured to apply pulse discharge so that a pulse current flows through the shunt resistor 7, and also to detect, for example, the intensity of a current having flowed through the resistor element 10 at a time of the pulse discharge. The state of a battery can be determined based on, for example, a current value detected at this time. Since a method for determining the battery state through pulse discharge is known, a detailed description thereof is omitted.\nConventionally, the battery state detection apparatus has adopted the known four-terminal method for the detection of a current. As well known, adoption of the four-terminal method enables accurate measurement of a current. FIG. 10(b) shows an equivalent circuit that detects a current having flowed through the resistor element 10 by means of the four-terminal method.\nAs shown in FIG. 9, the conventional board connection terminal 15 is configured to have a slit formed therein so that two connecting parts 30 are provided. More specifically, the board connection terminal 15 arranged at the first conductor part 11 side is provided with two connecting parts 30c and 30d. On the other hand, the board connection terminal 15 arranged at the second conductor part 12 side is provided with two connecting parts 30a and 30b. Thus, four connecting part 30 in total are arranged on the shunt resistor 7, which enables a current having flowed through the resistor element 10 to be measured by using the four-terminal method.\nIn the conventional battery state detection apparatus, to achieve the equivalent circuit shown in FIG. 10(b), the four connecting parts 30 (30a, 30b, 30c, 30d) are arranged in a straight line (FIG. 10(a))."}
{"text": "1. Field\nExemplary embodiments of the present disclosure relate to an ultrasound probe for generating images of the inside of a subject using ultrasonic waves.\n2. Description of the Related Art\nAn ultrasound diagnostic system includes a noninvasive apparatus that irradiates an ultrasound signal to a body surface at a target organ in the body and obtains cross-sectional images of, for example, soft tissue and blood flow.\nCompared to other imaging diagnostic systems such as X-ray diagnostic systems, computed tomography (CT) scanners, magnetic resonance imaging (MRI) systems and diagnostic systems for nuclear medicine, the ultrasound diagnostic system may have a compact size and low price, display images in real time, and provide a high level of safety by eliminating exposure to radiation. For at least these reasons, an ultrasound diagnostic system has been widely used for diagnosis in, for example, cardiac medicine, abdominal imaging, urology, obstetrics and gynecology.\nThe ultrasound diagnostic system includes an ultrasound probe which transmits an ultrasound signal to a subject and receives an ultrasound echo signal reflected from the subject to obtain an ultrasound image of the subject.\nThe ultrasound probe includes a piezoelectric layer of piezoelectric materials which convert electric signals into sound signals (e.g., acoustic), and vice versa, through vibration of the piezoelectric materials, a matching layer to reduce a difference in acoustic impedance between the piezoelectric layer and the subject to allow ultrasonic waves transmitted from the piezoelectric layer to be transferred to the subject as much as possible, and a lens to focus the ultrasonic waves traveling from the front of the piezoelectric layer on a specific point, and a backing layer to block the ultrasonic waves from traveling in an opposite direction from the rear of the piezoelectric layer to prevent image distortion."}
{"text": "1. Field of the Invention\nThis invention relates to a faucet assembly, more particularly to a faucet assembly with a valve which controls flow of water in a water mixing chamber to one of two outflow ports.\n2. Description of the Related Art\nReferring to FIGS. 1 and 2, a conventional faucet assembly is shown to include a faucet body 10 which has a water mixing chamber 103 for introducing hot and cold water from hot and cold water inflow passages 101,102. A primary valve 11 is disposed in the water mixing chamber 103, and extends in an axial direction to connect with a handle 12 so as to be rotatable to control the flow of the mixed water formed in the water mixing chamber 103 into an accommodating cavity 104 which is in fluid communication with the water mixing chamber 103. The accommodating cavity 104 is further in fluid communication with first and second outlets 105,106, such as a faucet spout and a showerhead. A flow direction controlling body 108 is disposed downstream of the water mixing chamber 103 and upstream of the second outlet 106. The flow direction controlling body 108 has an annular wall 107 which extends in the axial direction such that a secondary valve 13 is disposed sealingly and rotatably in the flow direction controlling body 108 by means of seal rings 131. A handle 14 has a spline connection with an outer portion 133 of the valve 13 to actuate rotation of the secondary valve 13. The valve 13 has a through hole 132 which extends radially and which is rotatable to communicate the accommodating cavity 104 with the second outlet 106 so as to permit the mixed water to flow out from the second outlet 106. When the through hole 132 is rotated so as not to register with the second outlet 106, the mixed water will flow out from the first outlet 105.\nHowever, a clearance arises between the outer surrounding wall of the valve 13 and the inner peripheral wall of the flow direction controlling body 22, thereby resulting in dripping at the second outlet 106.\nThe object of the present invention is to provide a faucet assembly which can overcome the aforementioned problem commonly associated with the prior art.\nAccording to this invention, the faucet assembly comprises a water mixing chamber which includes first front and rear inner walls facing towards and spaced apart from each other in an axial direction to define an accommodating space for receiving water. The first rear inner wall has hot and cold water inlets which are disposed to introduce hot and cold water flows into the accommodating space, and which can communicate with hot and cold water inflow passages, respectively. A transit chamber is in fluid communication with and downstream of the water mixing chamber. A valve assembly is disposed downstream of the mixing chamber and upstream of the transit chamber, and includes a primary valve with first and second valve portions. The first valve portion is disposed in the water mixing chamber, and is rotatable relative to the first rear inner wall about a first axis oriented in the axial direction to control the flows of hot and cold waters into the water mixing chamber via the hot and cold water inlets, respectively. The second valve portion is disposed between the water mixing chamber and the transit chamber to control the flow of mixed water formed in the water mixing chamber into the transit chamber. The valve assembly further includes a regulating member which is disposed in the water mixing chamber, and which is coupled to be rotatable with the first valve portion, and to be associated with the second valve portion. The regulating member extends in the axial direction and outwardly of the first front inner wall so as to be externally actuated to rotate the first valve portion and to control the second valve portion. A flow direction controlling body includes second front and rear inner walls which face toward and which are spaced apart from each other in the axial direction, and an inner tubular wall which surrounds a second axis oriented in the axial direction, which is interposed between, and which cooperates with the second front and rear inner walls to confine an accommodation cavity. The inner tubular wall has a communicating port which is disposed downstream of and which is communicated with the transit chamber, and which extends in a radial direction relative to the second axis to communicate the transit chamber with the accommodation cavity. The second rear inner wall has first and second outflow ports which are spaced apart from each other and which can respectively communicate with first and second outlets. A secondary valve includes a bottom wall with an annular periphery. The bottom wall is disposed to face toward and is in sliding contact with and rotatable relative to the second rear inner wall between first and second positions. The secondary valve further includes an annular wall which extends from the annular periphery in the axial direction towards the second front inner wall. The annular wall includes proximate and distal portions relative to the bottom wall. The distal portion is in sealing and sliding contact with the inner tubular wall when the bottom wall is rotated relative to the second rear inner wall. The proximate portion has a first passage extending in a radial direction relative to the second axis and communicating with the communicating port. The bottom wall has a second passage which extends therethrough in the axial direction to communicate with the first passage and which further extends angularly about the second axis. As such, when the bottom wall is rotated to the first position, the second passage communicates with the first outflow port. When the bottom wall is moved to the second position, the second passage communicates with the second outflow port. The secondary valve further includes a stem portion which is secured to the distal portion and which extends in the axial direction and outwardly of the second front inner wall so as to be actuated externally to rotate-the-bottom wall."}
{"text": "Presently there are in use ski boots constituted by a shell to which there is downwardly connected a sole, having a tip and a heel connectable to predisposed bindings fixed to a ski; the shell, usually provided in rigid or semi-rigid plastic material, covers approximately entirely the foot of the user.\nTo the shell there is rotatably connected, typically at the malleoli area, a cuff, adapted to partially cover the leg of the user to provide a support thereto during skiing.\nInternally to the shell and to the cuff there is typically positioned an inner boot, provided in soft material, adapted to increase the comfort-fit for the user.\nThe main drawback of such known ski boots is that they must be produced, for being adequate to the size of the foot of the user, in various dimensions: this increases production costs, there being necessary different machines and molds for producing the components in the different dimensions, and moreover there exists the problem of stocking the various sizes, for both the producer and for the seller.\nAs a possible solution to such drawback ski boots are usually produced in a number of limited sizes so as to be able to be used, each one, by users having foot sizes which are slightly different.\nThis however forces the user to buy a size which does not correspond to the exact dimensions of his own foot, forcing him or her to employ other expedients, such as for example the use of socks of increased thickness, for increasing comfort and for allowing thus to transmit in an optimum manner the forces from the foot to the ski.\nSuch drawback is greater in the case of users which are children, whose foot typically has dimensions which increase in relatively short times.\nAnother problem which should be held in mind is the fact that to a fixed size there corresponds a fixed shape of the calf of the user, for which the use of larger sizes and the compensation of the greater space for the foot with a sock may bring about localized areas of compression of the upper end of the cuff on the calf."}
{"text": "Previously, many types of lids or covers have been used in endeavoring to provide an effective means to enclose a container in a removable manner.\nThe prior art listed below did not disclose patents that possess any of the novelty of the instant invention; however the following U.S. patents are considered related:\nU.S. Pat. No.InventorIssue Date5,433,339SarverJul. 18, 19956,202,877 B1La Torre et al.Mar. 20, 20016,662,978 B2Lin et al.Dec. 16, 20036,913,159 B1GoldbergJul. 5, 2005D 509,408 SWard et al.Sep. 13, 20057,021,486 B1HurlbutApr. 4, 20067,416,093 B2Lin et al.Aug. 26, 2008D 588,002 SD'AmatoMar. 10, 2009D 592,913 SPinelli et al.May 26, 2009\nSarver in U.S. Pat. No. 5,433,339 teaches a cylindrical drinking vessel having a pair of integrally formed circumferential grooves and a pair of vertically oriented parallel handle support braces on its outer surface. A carabiner, or rope connection latched oval or a D-shaped ring fits between the handle support braces in a position to straddle the circumferential grooves. A pair of flexible bands are releasably fastened though the carabiner and around the vessel to secure the carabiner as a handle.\nU.S. Pat. No. 6,202,877 B1 issued to La Torre et al. provides a lip-openable lid for use in a container having a valve seat and an opening for egress of container contents. A container is also provided to employ the lid which has a gasket with an annular wall, a central portion radially inward of the wall and a peripheral flange engagable with the valve seat. The lid includes a rigid support that has liquid flow openings which supports the central portion of the gasket and positions the gasket against the valve seat.\nLin et al. in U.S. Pat. No. 6,662,978 B2 discloses a liquid storage vessel stopper having interchangeable plugs that utilizes stopper body with an inner shell affixed to an outer shell. The inner shell is configured to communicate with a liquid storage vessel interior and the outer shell is profiled to be threadably attached to the vessels outside surface. An interchangeable stopper plug is either a twist to pour plug or a push to pour plug type, the former has a main core smaller than the inner shell and incorporates a lid creating a dead air space forming an insulating barrier. When the main core is manually rotated at least a quarter of a turn the plug is unseated and a flow path is opened.\nGoldberg in U.S. Pat. No. 6,913,159 B1 teaches a water bottle strap that provides a manual engagement for bottles. A handle is defined from a substrate defining holes through which fingers or other digits may pass. The handle is held in tension along the side of the bottle by a top loop engaging the neck of the bottle and a base loop engaging the base of the bottle. The top loop and base loop are connected to the handle by means of an upper and lower connecting segment. Each strap is made of flexible and stretchable material so as to snugly engage the bottle in tension.\nU.S. Pat. No. 7,021,486 B1 issued to Hurlbut is for a drinking flask cap having a three-position push-pull lid. The top of the lid forms a drinking basis and the drinking basis can be sealed with a hygienic cap. The cap can be removed and placed on the bottom of the flask of storage and provides a stable base for the flask. A bottle portion of the flask is preferably made form a molded plastic material and can be provided with a finger hold such that the flask can be drunk from in a manner similar to a coffee cup.\nU.S. Pat. No. 7,416,093 B2 in U.S. patent Lin et al. discloses a container having a controlled opening lid that includes a container body with an open end providing access to an interior space configured to receive user supplied contents and a lid member configured to cover the open end of the container body when in a first position and permit access to the interior space when in a second position. A hinge is configured to move the lid member from the first position to the second position, the hinge including at least one dampening mechanism for damping a biasing force provided by the hinge when the lid member moves from the first position to the second position.\nFor background purposes and as indicative of the art to which the invention is related reference may be made to the remaining cited design patents issued to Ward et al. in U.S. Pat. D509,408 S, D'Amato in U.S. Pat. D588,002 S and Pinelli et al. in U.S. Pat. D592,913 S."}
{"text": "In recent years, from the viewpoint of protecting an occupant, the safety standard against the side collision of an automobile have been augmented, and the automotive door is required to exert the side collision performance sufficiently. In order to protect the occupant from the side collision, it is important to reduce the intrusion quantity of a colliding vehicle into the door part and to inhibit deformation of the door toward inside the cabin. In the side collision, the space for absorbing the collision energy is small compared with that in the collision of the front side and rear side. Against such collision, conventionally, a door beam was disposed inside the door along the vehicle traveling direction and deformation of the door toward inside the cabin was inhibited.\nAt present, a door beam of high tensile steel pipe and pressed material (Patent Document 1) or an aluminum extruded shape (Patent Document 2) has been developed and used. In general, for these door beams, one with a linear shape without a curve is used. In the Patent Document 3, a bending work method is described wherein a dent is not caused on an outside wall in bending in bending work of an extruded shape. Also, in the Patent Document 4, an embodiment is described wherein a door beam of a hollow body having a partitioning wall with a predetermined range being quenched is arranged in a direction oblique to a door.    Patent Document 1: Gazette of the Japanese Unexamined Patent Publication No. H8-216684    Patent Document 2: Gazette of the Japanese Unexamined Patent Publication No. H11-264044    Patent Document 3: Gazette of the Japanese Unexamined Patent Publication No. 2001-71038    Patent Document 4: Gazette of the Japanese Unexamined Patent Publication No. 2001-287608"}
{"text": "1. Field of the Invention\nThe invention relates to concepts associated with energy management systems and, more particularly, to means for providing remote communications associated with management of functions such as load balancing, load shedding and the like.\n2. Background Art\nA significant amount of work is currently being performed in technologies associated with control of what can be characterized as “environmental” systems. Such systems may be utilized in commercial and industrial buildings, residential facilities, transportation systems and other environments. Control functions may vary from relatively conventional HVAC temperature control to extremely sophisticated systems for control of the entirety of a city's subway complex.\nDevelopment is also being undertaken in the field of network technologies for controlling environmental systems. References are often made in current literature to “smart” buildings or rooms having automated and centralized environmental functionality. This technology provides for networks controlling a number of separate and independent functions, including temperature, lighting and the like.\nThere are a number of issued patents directed to various aspects of control of environmental systems. For example, Callahan, U.S. Pat. No. 6,211,627 B1 issued Apr. 3, 2001 discloses lighting systems specifically directed to entertainment and architectural applications. The Callahan lighting systems include apparatus which provide for distribution of electrical power to a series of branch circuits, with the apparatus being reconfigurable so as to place the circuits in a dimmed or “not-dimmed” state, as well as a single or multi-phase state. Callahan further discloses the concept of encoding data in a form detectable in electrical load wiring and at the load. The data may include dimmer identification, assigned control channels, descriptive load information and remote control functionality. For certain functions, Callahan also discloses the use of a handheld decoder.\nD'Aleo et al., U.S. Pat. No. 5,191,265 issued Mar. 2, 1993 disclose a wall-mounted lighting control system. The system may include a master control module, slave modules and remote control units. The system is programmable and modular so that a number of different lighting zones may be accommodated. D'Aleo et al. also disclose system capability of communicating with a remote “power booster” for purposes of controlling heavy loads.\nDushane et al., U.S. Pat. No. 6,196,467 B1 issued Mar. 6, 2001 disclose a wireless programmable thermostat mobile unit for controlling heating and cooling devices for separate occupation zones. Wireless transmission of program instructions is disclosed as occurring by sonic or IR communication.\nOther patent references disclose various other concepts and apparatus associated with control systems in general, including use of handheld or other remote control devices. For example, Zook et al., U.S. Pat. No. 4,850,009 issued Jul. 18, 1989 disclose the use of a portable handheld terminal having optical barcode reader apparatus utilizing binary imaging sensing and an RF transceiver. Sheffer et al., U.S. Pat. No. 5,131,019 issued Jul. 14, 1992 disclose a system for interfacing an alarm reporting device with a cellular radio transceiver. Circuitry is provided for matching the format of the radio transceiver to that of the alarm reporting unit. Dolin, Jr. et al., U.S. Pat. No. 6,182,130 B1 issued Jan. 30, 2001 disclose specific apparatus and methods for communicating information in a network system. Network variables are employed for accomplishing the communication, and allow for standardized communication of data between programmable nodes. Connections are defined between nodes for facilitating communication, and for determining addressing information to allow for addressing of messages, including updates to values of network variables. Dolin, Jr. et al., U.S. Pat. No. 6,353,861 B1 issued Mar. 5, 2002 disclose apparatus and methods for a programming interface providing for events scheduling, variable declarations allowing for configuration of declaration parameters and handling of I/O objects.\nAlthough a number of the foregoing references describe complex programming and hardware structures for various types of environmental control systems, it is desirable for certain functions associated with environmental control to be readily useable by the layperson. This is particularly true at a specific location, where it may be desirable to readily initially configure or reconfigure relationships or “correlation” between, for example, switching devices and lighting apparatus. Also, it may be desirable for such capability of initial configuration or reconfiguration to preferably occur within the proximity of the switching and lighting apparatus, rather than at a centralized or other remote location.\nIn addition to the foregoing, particular attention is being directed to energy conservation. In this regard, reviews are currently being undertaken with respect to time periods and quantities of power which are used in various types of residential, commercial and industrial facilities. For example, it would be desirable to be able to determine target values for power consumption within a total facility, and be able to adjust energy transmissions based on an “as needed” basis, so as to maintain power consumption at a specific level, or otherwise within a specific tolerance window."}
{"text": "The present invention relates to a new and distinctive soybean variety designated XB004C12, which has been the result of years of careful breeding and selection in a comprehensive soybean breeding program. There are numerous steps in the development of any novel, desirable soybean variety. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The breeder's goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include, but are not limited to: higher seed yield, resistance to diseases and/or insects, tolerance to drought and/or heat, altered fatty acid profile(s), abiotic stress tolerance, improvements in compositional traits, and better agronomic characteristics.\nThese product development processes, which lead to the final step of marketing and distribution, can take from six to twelve years from the time the first cross is made until the finished seed is delivered to the farmer for planting. Therefore, development of new varieties is a time-consuming process that requires precise planning, efficient use of resources, and a minimum of changes in direction.\nSoybean (Glycine max) is an important and valuable field crop. Thus, a continuing goal of soybean breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties.\nThe soybean is the world's leading source of vegetable oil and protein meal. The oil extracted from soybeans is used for cooking oil, margarine, and salad dressings. Soybean oil is composed of saturated, monounsaturated, and polyunsaturated fatty acids. It has a typical composition of 11% palmitic, 4% stearic, 25% oleic, 50% linoleic, and 9% linolenic fatty acid content (“Economic Implications of Modified Soybean Traits Summary Report”, Iowa Soybean Promotion Board & American Soybean Association Special Report 92S, May 1990). Changes in fatty acid composition for improved oxidative stability and nutrition are also important traits.\nIndustrial uses for processed soybean oil include ingredients for paints, plastics, fibers, detergents, cosmetics, and lubricants. Soybean oil may be split, inter-esterified, sulfurized, epoxidized, polymerized, ethoxylated, or cleaved. Designing and producing soybean oil derivatives with improved functionality, oliochemistry, is a rapidly growing field. The typical mixture of triglycerides is usually split and separated into pure fatty acids, which are then combined with petroleum-derived alcohols or acids, nitrogen, sulfonates, chlorine, or with fatty alcohols derived from fats and oils.\nSoybean is also used as a food source for both animals and humans. Soybean is widely used as a source of protein for animal feeds for poultry, swine, and cattle. During processing of whole soybeans, the fibrous hull is removed and the oil is extracted. The remaining soybean meal is a combination of carbohydrates and approximately 50% protein.\nFor human consumption soybean meal is made into soybean flour which is processed to protein concentrates used for meat extenders or specialty pet foods. Production of edible protein ingredients from soybean offers a healthy, less expensive replacement for animal protein in meats and dairy products."}
{"text": "Barcode identification systems having a number of barcode printers are known for printing jobs that include barcodes and associated alphanumeric characters on a web of record members such as tags, labels and the like. Typically barcode printers can print jobs in a limited number of formats. The format of a tag or label defines the type of data field, i.e., barcode fields, alpha fields, numeric fields, etc.; the position of the fields on the tag/label; the font used to print alphanumeric fields if more than one font is available; magnification and rotation factors if any; etc. as is well known to those in the art of barcode identification systems. Printers of various types typically have different format options. Because of the differences in printer format options, a job having a format capable of printing on one printer may not be able to print on a printer of a different type.\nTypical barcode printer systems having a number of printers, typically require that the printers of the system be of the same family such that the printers are compatible in that they communicate in the same language and are capable of printing data in the same formats. One known barcode information system utilizes a number of incompatible barcode printers as well as compatible barcode printers. As used herein, incompatible printers refers to printers that do not communicate in the same language and/or do not have substantially the same format options. In this known barcode identification system, a different queue is provided for each type of printer in the system. Although incompatible printers cannot be assigned to the same queue, each queue can have assigned thereto a number of printers of the same family, i.e., compatible printers. In this system, formats are manually created by entering format information without being prompted to by the system. Each format is created specifically for a queue so that it will be valid for the printers attached thereto.\nAnother known barcode printer will prompt a user to create a format for a single printer type. This system however cannot prompt a user to create a format for multiple printers of different types, i.e., printers that are incompatible. Therefore, if a job is to be printed on a number of different incompatible printers, different formats must be created for each of those printers in order for the printers to print the same job."}
{"text": "Transponders are electronic devices incorporated into secure documents such as “smart cards” and “electronic passports” using RFID (radio frequency identification) technology.\nThe transponder (or inlay, or chip card) itself generally comprises (includes):                a substrate (or inlay substrate) which may comprise a sheet of a synthetic material;        a chip (or chip module, or chip unit) installed on the substrate (or in a recess in the surface of the substrate) and having terminals (or contact surfaces, or pads); and        an antenna wire (or conductor) mounted on the substrate, formed with “turns” as a flat coil and connected (bonded) by its ends (see FIG. 2A) or end portions (see FIG. 2B) to the terminals of the chip. (In some of the drawings presented herein, only one end or end portion of the antenna wire may be shown, for illustrative clarity, particularly in the cross-sectional views.)        \nThe following dimensions and materials are merely exemplary:                the chip module may be generally rectangular, measuring for example 5 mm×8 mm (width×height), and may have a thickness of approximately 0.5-1 mm                    a first terminal disposed on one end of the chip module may measure for example 4 mm wide×2 mm high—in other words, nearly the entire width of the chip module, and approximately ¼-⅓ its height.            a second terminal disposed on one the other end of the chip module may measure for example 4 mm wide×2 mm high—in other words, nearly the entire width of the chip module, and approximately ¼-⅓ its height.                        The antenna wire may be “heavy” wire (such as 60 μm in diameter) requiring higher bonding loads than those used for “fine” wire (such as 30 μm). Rectangular section copper ribbon (such as 60 μm×30 μm) can be used in place of round wire. The antenna wire can be self-bonding copper wire or partially coated self-bonding copper wire, enamel copper wire or partially coated enamel wire, silver coated copper wire, un-insulated wire, aluminum wire, doped copper wire or litz wire.        The inlay substrate may be approximately 10 thousandths of an inch (300 microns) thick and may be a “synthetic paper” material such as Teslin™. The inlay substrate may have an area designated as a “transponder site” whereat the chip module will be installed. (A recess in the inlay substrate may constitute the transponder site.) The transponder site may itself have two areas designated as “terminals areas” corresponding in position to the two terminals of the chip module which will be installed at the transponder site. (The transponder site and terminal areas are generally geometric abstractions, the chip module and terminals are physical elements.) Hence, it should be understood that, where applicable, the terms (and reference numerals for) “transponder site” and “chip module” may be used interchangeably, and that the terms “terminal areas” and “terminals” may similarly be used interchangeably.        the chip module may be disposed in a recess extending into the surface of the substrate measuring for example 5.5 mm wide×8.5 mm high (generally the recess is only slightly larger than the chip module to allow some clearance during installation, while maintaining good registration)        \nA conventional method to produce an inlay site on an inlay substrate containing a high frequency RFID chip (or chip module) and an antenna embedded into a multi-layer substrate and connected to the terminals (terminal areas) of the RFID chip is disclosed in U.S. Pat. No. 6,233,818 and comprises                first positioning the RFID chip in a recess, supported by a lower substrate layer (of the multi-layer substrate), then start embedding (countersinking) a wire conductor onto or into the top substrate layer in the direction of the RFID chip,        then guiding the wire conductor over a first terminal area of the RFID chip, then continuing the embedding process by forming an antenna in the top substrate layer with a given number of turns,        then guiding the wire conductor over the second terminal area, and        finally embedding the wire conductor again into the top substrate layer before cutting the wire to complete the high frequency transponder site.        \nIn a next stage of the production process, the wire ends passing over the terminal areas are interconnected by means of thermal compression bonding. Adhesively placing a wire conductor onto the top substrate layer is an alternative to embedding, and typically involves self-bonding coated wire conductor.\nAn Inlay and Transponder of the Prior Art\nFIGS. 1A and 1B illustrate an inlay sheet 100 having a plurality of transponder areas (or sites). A selected one of the transponder areas 102 is shown in some detail. The vertical and horizontal dashed lines (in FIG. 1A) are intended to indicate that there may be additional transponder areas (and corresponding additional transponders) disposed to the left and right of, as well as above and below, the transponder area 102, on the inlay sheet 100. Typically, several transponders are fabricated on a single inlay sheet. In the main hereinafter the fabrication of a single transponder at a single transponder site will be discussed, except as may otherwise be noted.\nThe inlay sheet 100 may be a multi-layer substrate 104 comprising one or more upper (top) layers 104a and one or more lower (bottom) layers 104b. A recess 106 may be formed in (completely through) the upper layer 104a so that a chip module 108 may be disposed in the recess 106, and supported by the lower layer 104b. Alternatively, the substrate 104 may be a single layer substrate (not shown, but imaging layers 104a and 104b being a single layer rather than two distinct layers) with the recess extending only partially through the substrate.\nThe chip module 108 is shown having two terminals 108a and 108b on a top surface thereof. An antenna wire 110 is connected by its two end portions 110a and 110b to the corresponding two terminals 108a and 108b, respectively, of the chip module 108. In the main hereinafter the mounting and connection of only one end or end portion of an antenna wire to a corresponding one terminal of a chip module may be discussed, and may be taken as representative of how the other end or end portion of the antenna wire is connected to the other terminal of the chip module, except as may otherwise be noted.\nThe combination of chip module 108 and antenna wire 110 connected to the chip module at a transponder site 102 may be considered to be a transponder, and may be referred to by the same reference numeral 102. Additional layers of material (not shown) may be applied (laminated) to the transponder to make it suitable for use as a secure document such as an electronic passport or smart card. Thus, what is shown in FIGS. 1A and 1B can be considered to be an “interim product”. In the main hereinafter transponders which are interim products may be discussed except as may otherwise be noted.\nThe transponder 102 may be formed on an inlay substrate 104 having multiple layers comprising an upper layer 104a and a lower layer 104b, and may be exemplary of a smart card embodiment of a transponder for a secure document. A hot lamination process may be used to bond the upper and lower layers 104a and 104b together. The following dimensions are merely exemplary:                the overall thickness of the inlay substrate 104 may be approximately 450 μm        the thickness of the top layer 104a may be approximately 400 μm, and may comprise one or more layers of material        the depth of the recess 106 may be the same as thickness of the top layer(s) 104a         the thickness of the transponder chip 108 may be approximately 320 μm        the thickness of the bottom layer 104b may be approximately 240 μm, and may comprise one or more layers of material        \nGenerally, the recess 106 is sized and shaped to accurately position the chip module 108, having side dimensions only slightly larger than the chip module 108 to allow the chip module 108 to be located within the recess 106. And typically a dab of glue (not shown) in the bottom of the recess 106 will retain the chip module in the recess 106 during manufacturing. The following dimensions are merely exemplary:                the chip module 108 may measure approximately 5.0×8.0 mm        the recess 106 may measure approximately: 5.1×8.1 mm        the terminals 108a/b may measure approximately 5.0×1.45 mm        the wire may have a diameter between 60 and 112 μm        \nOne millimeter (mm) equals one thousand (1000) micrometers (μm, “micron”).\nIn FIGS. 1A and 1B, the recess 106 may be illustrated with an exaggerated gap between its inside edges and the outside edges of the chip module 108, for illustrative clarity. In reality, the gap may be only approximately 50 μm-100 μm (0.05 mm-0.1 mm).\nIn FIG. 1A the terminals 108a and 108b of the chip module 108 are shown reduced in size (narrower in width and or height), for illustrative clarity. (From the dimensions given above, it is apparent that the terminals 108a and 108b can extend substantially the full width of the chip module 108.)\nIt should be understood that the chip module 108 is generally snugly received within the recess 106, with dimensions suitable that the chip module 108 will not move around after being located within the recess 106, in anticipation of end portions 110a and 110b of the antenna wire 110 being bonded to the corresponding terminals 108a and 108b, respectively, of the chip module 108.\nAs best viewed in FIG. 1A, an antenna wire 110 is disposed, such as by embedding, on a top surface (side) of the substrate 104, and may be formed into a flat (generally planar) coil, having a number (such as four or five of “turns” (three shown), two ends, two end portions 110a and 110b, and a main intermediate portion between the two end portions. The overall length of the antenna thus formed may be approximately 1 meter.\nAs best viewed in FIG. 1A, the end portions 110a and 110b of the antenna wire 110 are shown extending completely over the terminals 108a and 108b of the chip module 108, from one side of the chip module 108 to the other, without (yet) being connected to the terminals 108a and 108b of the chip module 108. The end portions 110a and 110b of the antenna wire 110 may subsequently be connected, such as by thermo-compression bonding, to the terminals 108a and 108b of the transponder chip 108, respectively.\nAs best viewed in FIG. 1A, an intermediate portion 110c of the antenna wire 110 which is that portion of the antenna wire 110 which is between the two end portions 110a and 110b, and comprising most of the antenna wire is mounted (as indicated by the symbols “X”) in the surface of the inlay substrate 102. Only a few of the X's along the length of the intermediate portion of the antenna wire are omitted, for illustrative clarity.\nAs best viewed in FIG. 1A, a number of points (“a”, “b”, etc.) are illustrated along the length of the antenna wire 110 are shown. The following terminology may be used to describe different portions of the antenna wire.                the antenna wire extends between the points “a” at one “end” of the antenna wire and a point “f” at another end of the antenna wire        a short (such as 1-3 mm) portion of the antenna wire between the points “a” and “b” may be referred to as a “end segment” of the antenna wire, and generally includes the “end”. Two end segments are shown,                    one between the points “a” and “b”, and            the other between the points “e” and “f”.                        another short (6-8 mm) portion of the antenna wire which may pass over the terminal of the chip module may be referred to as the “end portion” of the antenna wire. The end portions may also be referred to as “connection portions”. Two end portions (or connection portions) are shown,                    one between the points “b” and “c”, and            the other between the points “e” and “f”.                        a main or “intermediate portion” of the antenna wire is that portion of the antenna wire which is between the two end portions, such as between the points “c” and “d”, and typically may be approximately 1 meter in length, formed with a number (such as 4 or 5) turns.        \nAs best viewed in FIG. 1B, the antenna wire 110 is “mounted” to the substrate 104a, which may comprise “embedding” (countersinking) the antenna wire 110 into the surface of the inlay substrate 102, or “adhesively placing” (adhesively sticking) the antenna wire 110 on the surface of the inlay substrate 102. In the main, hereinafter, mounting the antenna wire by embedding is discussed. As a general proposition, the antenna wire gets mounted (or fixed) to the inlay substrate 102 using a sonotrode tool 116 which vibrates, feeds out the wire, and embeds it into the surface of the inlay substrate 102.\nExamples of embedding a wire in a substrate, in the form of a flat coil, and a tool for performing the embedding (and a discussion of bonding), may be found in U.S. Pat. No. 6,698,089 (refer, for example, to FIGS. 1, 2, 4, 5, 12 and 13 of the patent). See also FIGS. 1 and 2 of U.S. Pat. No. 6,233,818. Both of these patents are incorporated by reference herein.\nIt is known that a coated, self-bonding wire will stick to a synthetic (e.g., plastic) substrate because when vibrated sufficiently to soften (make sticky) the coating and the substrate.\nIn the case of Teslin (synthetic paper), a normal insulated wire would not properly embed into the material, it would detach. Therefore, one uses self-bonding wire which attaches to the material with a slight penetration of the wire in the material. To sink the wire into the material, it is generally necessary to pre-press the antenna with chip connected using a hot lamination press. In contrast thereto, when pre-forming a channel or groove for the wire (such as by using laser ablation or other techniques, mentioned below, see FIG. 14A, et seq.) it is not necessary to perform this pressing operation.)\nFIG. 1B shows the antenna wire 110 feeding out of a capillary 116 of an ultrasonic wire guide tool (or “sonotrode”). The capillary 116 is typically disposed perpendicular to the surface of the substrate 100. The capillary 116 is omitted from the view in FIG. 1A, for illustrative clarity. In the main hereinafter the terms “capillary” and “sonotrode” may be used interchangeably, unless may otherwise be noted.\nThe capillary 116 may be vibrated by an ultrasonic vibration mechanism, so that it vibrates in the vertical or longitudinal (z) direction, such as for embedding the wire in the surface of the inlay substrate, or in a horizontal or transverse (y) direction, such as for adhesively placing the wire on the surface of the substrate.\nIn FIG. 1B, the antenna wire 110 is shown slightly spaced (in drawing terminology, “exploded” away) from the substrate, rather than having been embedded (countersunk) in the surface of the substrate, for illustrative clarity.\nAs described in U.S. Pat. No. 6,233,818, in a first “mounting” stage the end portions of the antenna wire are guided over the terminals of the chip module, and in a second “connecting” stage the end portions are connected to the terminals using a thermode tool to essentially weld the end portion (which may be referred to as a connection portion) of the wire to the terminal. FIG. 1B shows “generic” bond head 118, poised to move down (see arrow) onto the end portion 110b of the antenna wire 110 to bond it to the terminal 108b. \nThe connecting (or interconnection) process can be inner lead bonding (diamond tool), thermo-compression bonding (thermode), ultrasonic bonding, laser bonding, soldering, ColdHeat soldering (Athalite) or conductive gluing.\nAs best viewed in FIG. 1A, due to the layout of the antenna coils (or turns), the antenna wire 110 may need to cross over itself. This is illustrated in the dashed-line circled area. (With a different layout of the antenna coils, such as in FIG. 4 of either U.S. Pat. No. 6,233,818 or 6,698,089, no such crossover is necessary.) In order to prevent shorting (electrical contact between different portions of the antenna coil) the antenna wire 110 should be an insulated wire, generally comprising a metallic core and an insulation (typically a polymer) coating. Also, the polymer coating facilitates the wire being “adhesively placed” on (stuck to) a plastic substrate layer (such as 104a).\nIn order to feed the wire conductor back and forth through the ultrasonic wire guide tool, a wire tension/push mechanism (not shown) can be used or by application of compressed air it is possible to regulate the forward and backward movement of the wire conductor by switching the air flow on and off which produces a condition similar to the Venturi effect. This technique may be used to create a sufficient length of residual wire for positioning adjacent a first terminal of a chip module at a transponder site. This is particularly relevant to the techniques disclosed for example in US 2010/0141453 (“brushing”).\nThe material for one or more layers of the inlay substrate may comprise Teslin™, a waterproof synthetic film, single-layer, uncoated with a thickness of 356 microns.\nTeslinA single layer of microporous, polyolefin-based, uncoated film that bonds readily and firmly with toners, inks, adhesives and laminating films. Teslin is a waterproof synthetic material that works well with an inkjet printer, laser printer, or thermal printer. Teslin is also single-layer, uncoated film, and extremely strong. In fact, the strength of the lamination peel of a Teslin sheet is 2-4 times stronger than other coated synthetic and coated papers. Teslin comes in the sizes of 7 mil to 18 mil, though only sizes 10 mil and 14 mil are sized at 8.5″ by 11″, for printing with most consumer printers. Also available are perforated versions of Teslin, specifically, 1up, 2up, 3up, 6up, and 8up. Teslin is used widely in theproduction of drivers licenses, voter ID cards, and other forms of identification card. Because of its commercial availability, Teslin is also often found used for counterfeit or fake IDs. “Teslin” is a registered trademark of PPG Industries, Inc. for synthetic printing sheet.Tyvek ™Tyvek is a brand of spunbonded olefin, a synthetic material made of high-density polyethylene fibers; the name is a registered trademark of the DuPont Company. The material is very strong; it is difficult to tear but can easily be cut with scissors or any other sharp object. Water vapor can pass through Tyvek, but not liquid water, so the material lends itself to a variety of applications: medical packaging, envelopes, car covers, air and water intrusion barriers (housewrap) under house siding, labels, wristbands, mycology, and graphics.\nThe material for one or more layers of the inlay substrate may comprise PVC, PC, PE, PET, PETE, TYVEK, TESLIN, Paper or Cotton/Noil. The inlay substrate can also have special markings such as luminous threads, water marks, microscopic filings and optical polymer memory for additional security.\nFIG. 1C shows a transponder site similar to that of FIG. 1A, with a difference being that the ends of the antenna wire rather than end portions thereof are connected to the terminals of the chip module.\nFIG. 1D shows a transponder site similar to that of FIG. 1B, with a difference being that the recess is a “pocket” type recess that extends only partially through the thickness of what is typically a single layer substrate. In contrast therewith, a recess which extends entirely through a substrate, such as the recess 106 shown in FIG. 1B extending entirely through the layer 104a, may be referred to as a “window” type recess.\nPortions of the Antenna Wire\nFIG. 1E shows an antenna wire, and various portions thereof, such as:                two ends which are generally the geometric “ends” of the elongate antenna wire        two end segments which are generally short portions of the antenna wire at each end of the antenna wire. The end segments include the ends, and may be connected to the terminals of the chip module, such as in US 2010/0141453 (“brushing”)        two end portions (or connection portions) which are also short portions of the antenna wire which may be connected to the terminals of the chip module, such as in U.S. Pat. No. 6,233,818 or U.S. Pat. No. 7,546,671 (“looping”)        a main or intermediate portion which is the longest portion of the antenna wire, between the two end portions, and which may be formed into several turns of a flat coil antennaSome Examples of “Final Products”        \nTransponders such as are shown in FIGS. 1A,B and 1C,D may be considered to be “interim products” in that some further steps or elements may be needed before getting the product into the “hands of the consumer”. For example, various cover layers may be laminated to the inlay substrate to protect (and secure) the transponder, as well as for imprinting with information. The end result, or “final product”, may be a secure document such as an electronic passport or a smart card.\nFIG. 1F shows an example of a security document which may be a National ID (identification) Card (or electronic ID, “eID” card) comprising a multi-layer (2 layer) inlay substrate, and additional layers comprising a top overlay layer and a bottom overlay layer. An RFID chip module and corresponding antenna (not shown) may be mounted in the inlay substrate(s). The chip module (not shown) may have a mold mass and a leadframe. The additional top and bottom layers may be anti-scratch layers, and protect the inlay substrate(s). The eID card, inlay substrate layer and top and bottom layers are not shown to scale.\nSome dimensions for and properties of the layers may be:\nTop overlay layertransparent 80 micronsInlay substrate moldwhite185 micronsInlay Substrate-Leadtransparent 80 micronsBottom Overlay Layertransparent 80 microns\nThe layers of the inlay substrate for a smart card may comprise PVC, which has limited life. Smart cards are often replaced (renewed) every few years.\nThe layers of the inlay substrate for a national ID care may comprise PC, which may be more durable (longer life) than PCV.\nFIGS. 1G and 1H illustrate an exemplary construction for an electronic passport cover, corresponding generally to the single layer inlay substrate construction shown in FIG. 1D. The inlay substrate for a US passport may comprise Teslin™.\nA cover layer may be disposed over the inlay substrate 1 for the final product. The material for the cover layer may be a cloth product, with chemistry in the coatings and a leather-like appearance to the cloth, such as by Holliston Inc. (905 Holliston Mills Road, Church Hill, Tenn. 37642; www.holliston.com)\nSome Examples of Chip Modules\nIn the main hereinafter, the discussion may focus on RFID chip modules which are leadframe-type modules. However, some of the techniques for producing security documents discussed herein may also be applicable to epoxy glass modules (chip on FR4, wire bonded, glob topped).\nFIG. 1I shows an example of an RFID chip module which is a “leadframe module” comprising:                a leadframe having a thickness of approximately 80 μm        an RFID chip disposed on and connected by wire bonds to the leadframe, having a thickness of approximately 80 μm        a mold mass disposed over the chip and wire bonds, having a thickness of approximately 240 μm        an antenna wire having end portions connected to “connection areas” of the leadframe, typically on a side of the leadframe opposite the RFID chip (as shown), but the end portions can also be connected to connection areas on the same side of the lead frame as the RFID chip.        \nThe total thickness of the leadframe module may be 320 μm, such as for an inlay substrate having a thickness of approximately 356 μm. Generally, the chip module will be disposed in a recess in the inlay substrate so as to be concealed therein.\nFIG. 1J shows an example of an RFID chip module which is an “epoxy glass module” comprising:                an interconnect substrate, such as FR4 (printed circuit board substrate material), having a thickness of approximately 100 μm (FR4 is 100 μm and the chip & glob top 160 μm=overall 260 μm)        an RFID chip, wire-bonded (alternatively flip-chip connected with solder bumps and underfiller, as illustrated) to the FR4 substrate, having a thickness of approximately 100 μm        a glob top epoxy disposed over the chip and connections, having a thickness with chip of approximately 160 μm        an antenna wire having ends connected to “connection pads”, typically on the same side of the FR4 substrate as the RFID chip, but can also be connected on the opposite side of the FR4 substrate as the chip.        \nThe total thickness of the epoxy glass module may be 260 μm, such as for an inlay substrate having a thickness of approximately 365 μm. Generally, the chip module will be disposed in a recess in the inlay substrate so as to be concealed therein.\nGenerally speaking, epoxy glass modules are inherently somewhat more flexible than leadframe modules. This is a factor that may need to be taken into consideration when incorporating an RFID module into a secure document. And, whereas leadframe modules are typically rectangular, the mold part (glob top) of an epoxy glass module are typically round.\nIt should be understood that, although FIG. 1J shows a flip chip connection between the RFID chip and the FR4 substrate, the chip can be wire-bonded to the substrate (such as was shown in FIG. 1I, for the leadframe-type module.)"}
{"text": "The present application relates to a method of forming a wiring of a light emitting device, a substrate for mounting a light emitting device, a display, a back light and an electronic appliance, particularly suitable for a substrate for mounting a light emitting device in which a minute light emitting device, for example, a light emitting diode is mounted on the substrate and then a wiring is formed thereon, and a display, a back light and an electronic appliance using the substrate for mounting a light emitting device.\nAs a light emitting device, for example, a light emitting diode, there are a light emitting diode in which a p-side electrode and an n-side electrode are formed on one side, and a light emitting diode in which a p-side electrode and an n-side electrode are separately formed on a light emission surface (the surface to take out light) and on a surface on the opposite side thereof (the back side). An example of the former light emitting diode is shown in FIGS. 14a and 14B, and an example of the latter light emitting diode is shown in FIGS. 15A and 15B. Here, FIGS. 14A and 14B show a plan view and a cross section, respectively. In addition, FIGS. 15A and 15B show a cross section and a bottom view, respectively.\nIn the light emitting diode shown in FIGS. 14A and 14B, the structure of a light emitting diode is formed of an n-type semiconductor layer 101, a light emitting layer (an active layer) 102 and a p-type semiconductor layer 103, in which a p-side electrode 104 is formed on the p-type semiconductor layer 103, and an n-side electrode 105 is formed on the n-type semiconductor layer 101. Connecting conductive materials 106 and 107 are formed on the p-side electrode 104 and the n-side electrode 105, respectively. In this case, the back side of the n-type semiconductor layer 101 is the light emission surface. A protective insulating resin 108 is formed so as to cover the surfaces except the light emission surface and the connecting conductive materials 106 and 107 in the light emitting diode.\nIn the light emitting diode shown in FIGS. 15A and 15B, the structure of a light emitting diode is formed of an n-type semiconductor layer 101, a light emitting layer 102 and a p-type semiconductor layer 103, in which a p-side electrode 104 is formed on the p-type semiconductor layer 103, and an n-side electrode 105 in a linear shape is formed on one side of the back side of the n-type semiconductor layer 101 that is the light emission surface. A connecting conductive material 106 is formed on the p-side electrode 104. A protective insulating resin 108 is formed so as to cover the surfaces except the light emission surface and the connecting conductive material 106 in the light emitting diode.\nFIGS. 16A and 16B show the state in which the light emitting diode shown in FIGS. 14A and 14B is mounted on a substrate and wirings are formed thereon. Here, FIGS. 16A and 16B show a plan view and a cross section, respectively. As shown in FIGS. 16A and 16B, in this example, the light emitting diode is mounted on a transparent substrate 201 as the n-type semiconductor layer 101 is turned down, an insulating resin 202 is buried so that the connecting conductive materials 106 and 107 are exposed around the light emitting diode, and then a wiring 203 which is connected to the connecting conductive material 106 and a wiring 204 which is connected to the connecting conductive material 107 are formed on the insulating resin 202. In this case, the wirings 203 and 204 are lead to the opposite sides to each other.\nFIGS. 17A and 17B show the state in which the light emitting diode shown in FIGS. 15A and 15B is mounted on a substrate and wirings are formed. Here, FIGS. 17A and 17B show a plan view and a cross section, respectively. As shown in FIGS. 17A and 17B, in this example, the light emitting diode is mounted on a wiring 203 which is formed in advance on a substrate 201 as the connecting conductive material 106 is turned down, an insulating resin 202 is buried so that the light emission surface is exposed around the light emitting diode, and then a wiring 204 which is connected to the n-side electrode 105 is formed on the insulating resin 202 as the wiring 204 is lead to the opposite side of the wiring 203.\nIn addition, in a light emitting diode array having a center electrode type light emitting diode part in which an ohmic contact layer is disposed at the center of a light emission surface and light emission surfaces on both sides thereof emit light, such a technique is proposed that the ohmic contact layer is configured of a base part which passes through the center of the light emission surface and branch parts which are extended in the direction crossing the direction in which a bonding electrode is led from the base part, and ohmic electrodes are provided on the ohmic contact layer as tracing the shapes of the base part and the branch part (see Patent Reference 1 (JP-A-2004-14676)). However, this technique is greatly different in technical principles from the light emitting device according an embodiment in which an electrode on the light emission surface is formed nearly in a linear shape having the width narrower than that of the light emission surface, a wiring connected to the electrode is formed nearly in a linear shape having the width narrower than that of the light emission surface, and the wiring is crossed to the electrode.\nFrom the viewpoint of easiness to form a wiring after mounted, as compared with the light emitting diode shown in FIGS. 14A and 14B in which the p-side electrode 104 and the n-side electrode 105 are formed on one side, the light emitting diode shown in FIGS. 15A and 15B is more preferable in which the n-side electrode 105 and the p-side electrode 104 are separately formed on the light emission surface and on the surface on the opposite side thereof. In the light emitting diode shown in FIGS. 15A and 15B, it is desirable that the n-side electrode 105 to be formed on the light emission surface is as small as possible in order to suppress the light blocking against the light emission surface due to the n-side electrode 105. In addition, in the case in which the light emitting diode is mounted and then the wiring 204 is formed which is connected to the n-side electrode 105 on the light emission surface, it is necessary to form the wiring 204 as well not to block light against the light emission surface as much as possible. Moreover, the light blocking area on the light emission surface is greatly changed in the cases in which the position is shifted in mounting the light emitting diode on the substrate 201 and in which the position is shifted in forming the wiring 204, causing the brightness is varied, which has to be prevented. However, in practice, it is actually difficult to satisfy all of these demands.\nThen, a method shown in FIGS. 18A and 18B is proposed by the inventor as a method of forming a wiring that satisfies these demands. Here, FIGS. 18A and 18B show a plan view and a cross section, respectively. As shown in FIGS. 18A and 18B, in the method, a transparent electrode 205 formed of ITO (indium tin oxide) is formed to cover the entire light emission surface so that an n-side electrode 105 formed on the back side of an n-type semiconductor layer 101 that is the light emission surface is connected to a wiring 204 formed on an insulating resin 202.\nHowever, in the method of forming a wiring shown in FIGS. 18A and 18B, in practice, it is difficult to connect the n-side electrode 105 on the light emission surface to the wiring 204 formed on the insulating resin 202 with no breaks of the transparent electrode 205, and materials for use in forming the transparent electrode 205 are relatively expensive, causing a problem that fabrication costs for a mounting substrate are increased."}
{"text": "There are a variety of devices the purpose of which is to evaporate the fuel components entering the intake in the liquid state by preheating the fuel-air mixture. In this manner, a better preparation of the fuel-air mixture is achieved, whereby the presence of constituents injurious to health in the exhaust may be diminished. Preheating of this type is advantageous especially at low load, since it prevents the carbureted fuel from recondensing in this operating mode, with resulting combustion of a leaner fuel-air mixture, thereby reducing the efficiency of the engine. At full load when the engine is hot, on the other hand, preheating is not desirable since losses of fill would result. To obtain controllable preheating, more or less effective depending on the operating condition of the engine, dampers have been proposed, arranged in the intake and/or exhaust line in order to deflect or divert the incoming fuel-air mixture and/or the exhaust flow from the heat transfer wall. Such dampers, however, have the disadvantage of interfering seriously with the direction of flow of the incoming fuel-air mixture by a deflection or partial constriction and displacement of the port area, and may lead to separation of the mixture, while in the exhaust line a resistance is created that undesirably increases the back pressure."}
{"text": "Monitoring and diagnosing the status of rotating and reciprocating machinery start with accurate and dependable measurements from a transducer and its associated electronics and then proceed to other sophisticated analyzing apparatus for reduction and display. One such transducer is a proximity transducer which may be utilized for, inter alia, monitoring the vibration characteristics of a rotating shaft of a machine. In this environment, the transducer must operate under very adverse physical, chemical and mechanical conditions and it is often very difficult to replace such transducers. Thus, there is an ongoing effort to make the proximity transducer one of the most reliable parts of the monitoring system.\nTypically, the proximity transducer, in conjunction with associated electronics, outputs a signal correlative to the spacing between an object or \"target\" (the rotating shaft of the machine) and a sensing coil of the proximity transducer. It is critical that the length or spacing between the target and the sensing coil of the proximity transducer remains within the linear range of the transducer for providing accurate and reliable measurements when in operation. Thus, the hallmark for providing accurate and reliable measurements relies on providing a transducer which is impervious to the predations of the environment and which does not consume an inordinate amount of the linear range of the transducer.\nThe following prior art reflects the state of the art of which applicant is aware and is included herewith to discharge applicant's acknowledged duty to disclose relevant prior art. It is stipulated, however, that none of these references teach singly nor render obvious when considered in any conceivable combination the nexus of the instant invention as disclosed in greater detail hereinafter and as particularly claimed.\n______________________________________ U.S. PATENT DOCUMENTS DOCUMENT NO. DATE NAME ______________________________________ 2,361,348 October 24, 1944 Dickson et al. 2,890,505 June 16, 1959 Brand 3,932,828 January 13, 1976 Plunkett et al. 4,000,877 January 4, 1977 Shead et al. 4,162,138 July 24, 1979 Byrne 4,377,548 March 22, 1983 Pierpont 4,408,159 October 4, 1983 Prox 4,419,646 December 6, 1983 Hermle 4,470,786 September 11, 1984 Sano et al. 4,680,543 July 14, 1987 Kohen 4,829,245 May 9, 1989 Echasseriau et al. 4,954,307 September 4, 1990 Yokoyama 4,959,000 September 25, 1990 Giza 5,016,343 May 21, 1991 Schutts 5,018,049 May 21, 1991 Mehnert 5,021,737 June 4, 1991 Schutts 5,039,942 August 13, 1991 Buchschmid, et al. 5,049,055 September 17, 1991 Yokoyama 5,122,046 June 16, 1992 Lavallee et al. 5,133,921 July 28, 1992 Yokoyama 5,138,292 August 11, 1992 Forster 5,147,657 September 15, 1992 Giza 5,151,277 September 29, 1992 Bernardon, et al. 5,182,032 January 26, 1993 Dickie et al. 5,226,221 July 13, 1993 Kilgore 5,240,397 August 31, 1993 Fay et al. 5,252,051 October 12, 1993 Miyamoto et al. 5,351,388 October 4, 1994 Van Den Berg, et al. 5,376,325 December 27, 1994 Ormson ______________________________________\n______________________________________ FOREIGN PATENT DOCUMENTS DOCUMENT NO. DATE COUNTRY ______________________________________ UK 1 313 748 April 18, 1973 Great Britain UK 1 353 603 May 22, 1974 Great Britain JA-139710 August 6, 1978 Japan WO 84/03794 September 27, 1984 PCT FR 2576-245-A July 25, 1986 France JA 6-37130-A October 2, 1994 Japan ______________________________________\nThe two patents to Schutts and the patent to Van Den Berg, et al. reflect assignee's ongoing commitment to providing an accurate sensor which is impervious to predations of the environment.\nThe French patent to Jaeger teaches the use of a method and apparatus for injection molding of an elongated detector with a sensor at one end. One end of the detector is supported by the mold while the sensor end is engaged by a centering sleeve (130). The centering sleeve (130) terminates in a piston (132) which is mobile in a cylinder (126) and slides about a fixed rod (138). Thermoplastic is injected into the mold and the centering sleeve is removed from the sensor when the injection process is only partially complete.\nThe Japanese patent to Kawakami teaches the use of sealing a semiconductor chip by transfer molding. A semiconductor chip (4) is mounted on a carrier (1) and is held at a fixed position via mobile pins (17) and (18). The mobile pins (17) and (18) are movably fitted to a top force (11) and a bottom force (12) for freely advancing or retracting the pins into and out of the cavities (15) and (16). While the pins (17) and (18) are in contact with the carrier (1), a resin (20) is injected into the cavities (15) and (16) through gates (13) and (14) and the pins (17) and (18) are gradually retracted in accordance with the injecting state of resin (20).\nThe patent to Yokoyama teaches the use of an apparatus for manufacturing plastic encapsulated electronic semiconductor devices. A support pad (11) is firmly fixed by first and second sliders (24A) and (24B) and also by mold halves (19) and (20). The first and second sliders (24A) and (24B) are moved outwardly of a cavity (25) when the cavity is half filled with the plastic encapsulating material. Spaces, which are formed after movement of the sliders, are filled with the plastic encapsulating material directly poured through a gate (23). It should be noted that the sliders (24A), (24B) perform two functions in the cavity (25). The first function is to firmly grasp a thin end portion of the support pad (11) thereby fixing same in position and the second function is to narrow a passage through the gate (23).\nThe patent to Pierpont teaches the use of a method for encapsulating an electrical component having radial leads. A large multi-cavity mold is loaded with a plurality of horizontal radially leaded capacitors. The closed mold grips the leads (15). A top pin (25) pushes each component body downward a predetermined degree so that each body is left in about the same position within the corresponding mold cavity. A bottom pin (28) then pushes each component body upwardly only slightly above a center of the mold cavity so that upon bottom pin withdrawal, stresses within the lead wires of the electrical component cause the component bodies to spring back just to the mold cavity center portion. Molding resin is then introduced.\nThe other prior art listed above, but not specifically discussed, teach other sensor devices and molding processes and further catalog the prior art of which the applicant is aware. These references diverge even more starkly from the references specifically distinguished above."}
{"text": "In today's networks, devices, such as switches or routers, are used not only to route and/or switch packets toward their destinations, but also to collect information related to network devices, applications that are hosted on the network devices, and/or the network traffic. In some instances, the devices may collect the information from one or more flows, each of which is a stream of packets that are forwarded from a source to a destination."}
{"text": "The efficiency of semiconductor and integrated circuit fabrication process is sensitive to not only the quality of the materials utilized in the fabrication procedure but also the environment of the fabrication location. In the prior U.S. Pat. No. 5,085,760, the invention was directed to sensing or monitoring the very low levels of oxygen, in the parts per billion range, present in the inert blanketing gases required in the fabrication procedures. A similar problem in the fabrication process is that the presence of very low levels of moisture in the gas feed stream can severely effect both the yield and the quality of the integrated circuits or the like. In the past, a variety of sophisticated analytical instruments to monitor the moisture in a gas feed stream from high parts per million (ppm) to sub-parts per billion (ppb) levels were developed. Many of these prior art instruments utilized for sub-parts per billion (ppb) level measurements are relatively complex and expensive. The prior art instruments based on the electrolytic mechanism for decomposing or splitting the water in gases are relatively simple and inexpensive but lack the sensitivity to measure moisture in sub parts per billion levels and exhibit the sluggish response and recovery time. These electrolytic hygrometers utilized the principles disclosed by F. A. Keidel in Analytical Chemistry, Vol. 12, page 2043(1959). This well known, prior art hygrometer consists of an electrolytic cell having a pair of electrodes covered with a hygroscopic electrolyte such as Phosphorous Pentoxide P.sub.2 O.sub.5. In this sensor, the water vapors in the gas are absorbed at the surface of the hygroscopic layer and electrolyzed to gaseous oxygen and hydrogen. The current drawn by the electrolytic cell is a direct measure of the water being electrolyzed. At equilibrium, the cell current and the gas flow rate gives the absolute and continuous measurement of the amount of water in the gas stream.\nElectrolytic cells with the different physical configurations based on the Keidel principles have been disclosed in the literature. The most noticeable and frequently used physical configuration is in the form of a tubular conduit having a pair of electrode wires helically positioned in parallel on the inner wall of the conduit from end to end or alternatively an interdigitated grid of electrodes deposited on an insulating surface that is coated with a thin layer of a hygroscopic film, P.sub.2 O.sub.5, derived from the dehydration of a thin layer of phosphoric acid by electrolysis at elevated temperatures. The major problems experienced with such a configuration is that the electrolytic layer(1) upon prolonged exposure to high moisture levels, is converted to phosphoric acid that tends to run off the electrodes and (2) upon the prolonged exposure to gases with moisture levels in the sub parts per billion level causes the electrolytic layer to crack. Both of these undesirable conditions cause discontinuity in the electrolyte layer and as a result the sensor loses its sensitivity. In addition, when the electrolyte layer is fully dehydrated, the sensor's response to small changes in the moisture content becomes very sluggish.\nA solid, perfluorinated, ion-exchange polymer of the prior art, commercially available from E.I. du Pont de Nemours & Company, Inc. of Wilmington, Del. and commercially identified as a \"Nafion\" element, has been used as the electrolyte element in an electrolytic sensor. This ion-exchange polymer in the form of a thin film sandwiched between a pair of electrodes for trace water sensing and is disclosed in the U.S. Pat. No. 4,514,278 granted on Apr. 30, 1985. U.S. Pat. No. 4,954,238 granted on Sep. 4, 1990, discloses a moisture sensing hygrometer element in the form of an ion-exchange film of the \"Nafion\" type cast on the pair of electrodes on a tubular substrate. These \"Nafion\" devices are considered to have less disadvantages than those inherent in the Keidel-phosphorous pentoxide films. However, these \"Nafion\" based humidity sensors are less sensitive by more than two orders of magnitude compared to those using P.sub.2 O.sub.5 with identical physical configurations and with a proportionate effect on the lower detection limit. The lower sensitivity with perfluorinated ion-exchange polymer electrolyte (\"Nafion\") is primarily due to the decreased ionic conductivity of the \"Nafion\" polymer in the presence of lower levels of moisture. It is well known that water is essential for solvating the hydrogen ions of the sulfonic acid groups of the \"Nafion\" polymer for the ionic conduction to occur in the \"Nafion\" membrane. This is considered to preclude the use of the ion-exchange polymer as the electrolyte element in either very dry environments or at temperatures above 100 degrees Centigrade with unhumidified gases and without the means to supply water to the polymer.\nIn the above referenced prior U.S. Pat. No. 5,164,053, an improved electrochemical sensing cell is disclosed that utilizes a solid, perfluorinated ion-exchange polymer (PFIEP) element as an electrolyte element in the sensor. The electrolyte element taught to the art by the aforementioned patent utilizes a solid polymer electrolyte that has been equilibrated with phosphoric acid for maintaining the ionic conductivity of the polymer up to temperatures on the order of 180 degrees Centigrade. This prior patent teaches the solvation of the hydrogen ions of the sulfonic acid groups of the ion-exchange by the phosphoric acid that provides with a hydrogen bonding network for maintaining the ionic conductivity of the polymer without the need to add water to the polymer. This general structure is represented as follows: ##STR1##\nIt has only recently become known that it has been reported in the literature that incorporating phosphoric acid in a solution of the \"Nafion\" material for an electrolytic moisture sensor has been found to function in relatively dry environments with better sensitivity than the known prior art. This is disclosed by Huang et al. in Analytical Chemistry, Vol. 63, No. 15 for Aug. 1, 1991 on pages 1570-1573. The known humidity sensors for detecting sub parts per million moisture levels are typically powered at 30-100 volts. Under such electrical operating conditions phosphoric acid is readily converted into the dehydrated form of phosphorous pentoxide, P.sub.2 O.sub.5. This causes the electrolyte layer to lose its continuity and thereby the sensor loses its sensitivity. I have found that in a phosphoric acid, equilibrated perfluorinated ion-exchange polymer electrolyte used for electrochemical sensors, where the voltage between two electrodes is typically less than 1.23 volts (1.23 volts is the thermodynic voltage between the hydrogen and oxygen electrodes), as disclosed in U.S. Pat. No. 5,164,053, that the phosphoric acid is not converted to the dehydrated form and thus the conductivity of the electrolyte element is maintained. I have also determined that a moisture sensor utilizing a phosphoric acid treated, solid, perfluorinated ion-exchange polymer electrolyte element where the voltage is typically between 30-100 volts, over a prolonged period of time and in a very dry environment, (moisture levels in the lower end of the parts per billion levels), phosphoric acid is partially converted to the P.sub.2 O.sub.5 form and thereby results in decreasing the sensitivity of the sensor.\nA further aspect of the prior art is found in the U.S. Pat. No. 4,900,405 granted on Feb. 13, 1990 for a \"Surface Type Microelectronic Gas and Vapor Sensor\". This prior art principally discloses various micro-sensing structures having two, three or more electrode structures deposited in close proximity to one another on the same side of an active area on a surface of the substrate so that the ion migration is fast. The substrates may be constructed of an electrically insulative material or of a semiconductor material. These structures utilize solid, electrolytes including ion-exchange polymers of the \"Nafion\" type and takes advantage of the hydrophobic nature of such a polymer. All of the disclosed structures, however, utilizing the \"Nafion\" type ion-exchange member require the use of an aqueous reservoir in contact with the solid electrolyte medium to keep it from drying out and inactivating the disclosed microsensor. This problem is solved by the disclosure in the aforementioned U.S. Pat. No. 5,164,053 and further improved upon herein.\nAnother piece of prior art, utilizing an aqueous electrolyte with 3 electrodes, discloses a means to enable the various impurities in the environmental air to be distinctly detected and measured is found in U.S. Pat. No. 3,776,832. The gases, noxious atmospheric pollutants were selected from the group consisting of carbon monoxide, nitric oxide, hydrocarbons, ethanol and methanol in air. This is accomplished by conveying the same air sample through a series of individual cells, with each cell being constructed and defined to detect only a single impurity. This prior art patent discloses the use of potentiostat circuit means for maintaining a constant or fixed relative potential difference between the anode or sensing electrode and the reference electrode of the 3 electrode sensors with the voltage selected in accordance with the gas species to be sensed.\nThere is, then, a present need for an improved, relatively inexpensive electrochemical/electrolytic sensor with a greater stability and faster responses and recovery times than those presently known and in commercial use that are not subject to drying out during continuous use."}
{"text": "1. Field of the Invention\nThe invention relates to instruments for measuring environmental evaporativity. Specifically, the invention relates to a device which models the diffusion resistance and albedo of field crops.\n2. Related Art\nA variety of atmometers or evaporimeters are known for measuring evaporativity of liquids. An atmometer is an instrument for measuring the evaporative capacity of air. The influence of the environmental complex on the vaporitization of water is known as evaporativity.\nThe simplest type of atmometer is an open pan. Water is placed in the pan, and a variety of methods are employed for determining the amount of water evaporation. Volume readings may be computed from depth measurements, or the pans may be weighed periodically in order to determine how much water has evaporated. Open pans have a number of disadvantages when used to estimate the crop water use or \"evapotranspiration\" of field crops. Pans are a free water surface which respond to heat absorption (advective energy) i.e., heat carried by the wind different than a crop surface. A free water surface does not exert any degree of resistance to the passage of water vapor into the air as a crop surface does. Also the large volume of water in a typical pan stores heat during the day resulting in evaporation at night which doesn't occur with a crop. Furthermore, wind tends to cause splashing out of the water, so that water is lost as a result of winds and not evaporation. In addition, the creation of ripples and waves increases the evaporating surface in an uncertain and variable manner.\nA second type of atmometer is a wet paper or cloth atmometer. This form of atmometer employs a paper or cloth evaporating surface which is kept wet automatically by means of a distilled water supply from a reservoir. The evaporating surface is often in the shape of a disk or circle and can be made from thick filter paper or blotting paper. Evaporation occurs around the cylindrical edge of the circle as weel as from the two planar surfaces. One difficulty encountered by wet paper and cloth atmometers is that they are necessarily exposed unequally to wind from different directions. These atmometers are also free water surfaces offering no diffusion resistance to the water vapor evaporated.\nA third type of atmometer is a porous porcelain atmometer. The porous porcelain atmometer comprises a porous porcelain piece having a hollow porous porcelain sphere or cylinder and a downwardly extending cylindrical neck. A rubber stopper seals the lower open end of the cylindrical neck. A tube extends from the cylindrical neck through the rubber stopper into a liquid reservoir. The porous porcelain piece is filled with distilled water and is connected to the water reservoir through the tube. Atmospheric pressure on the water in the reservoir keeps the tube and porcelain piece filled while evaporation takes place from the external porous porcelain surface. Small water-air menisci form over the external pore openings of the porcelain wall. This prevents air from entering into the porcelain structure. Evaporation takes place from the outer surface of the cylindrical neck and liquid water moves up from the reservoir and out through the porous walls at a rate adequate to keep the instrument filled and to maintain the microscopic menisci at the external pore openings of the porcelain structure.\nOne important application of atmometers is the estimation of water transpired by a green crop in a field. In measuring the crop consumptive water use or \"evapotranspiration,\" it is important to consider a number of factors which will influence this parameter. Two key factors are the diffusion resistance of the crop and the albedo (amount of incident radiation which is reflected) of the crop. The known atmometers, being free water surface evaporimeters, do not account for these factors, and therefore, the estimated crop consumptive evapotranspiration is inaccurate.\nThe three above-mentioned atmometers are described in Livingston, \"Atmometers of Porous Porcelain and Paper, Their Use in Physiological Ecology,\" Ecology, Vol. XVI, No. 3, July 1935.\nU.S. Pat. No. 3,540,277 discloses an evaporimeter which raises the level of liquid to compensate for decreases in liquid level due to evaporation. While the actual amount of evaporated liquid can be fairly accurately measured by the disclosed apparatus, such measurements will not render extremely accurate estimates of crop evapotranspiration, as the diffusion resistance and albedo of crops have not been taken into account in the disclosed apparatus.\nGerman Pat. No. 58953 discloses an evaporimeter having a valve at an upper end. Heat transmitted to the water in the evaporimeter expands and opens the upper valve to permit the evaporation of water. Here, the evaporation takes place from a free water surface, i.e., the surface is open to the atmosphere and does not account for diffusion resistance or albedo.\nU.S. Pat. No. 3,898,872 and U.S. Pat. No. 4,412,477 also disclose apparatus which permit water movement or water evaporation from a \"free\" surface, which, as mentioned above, neglect diffusion resistance and albedo.\nU.S. Pat. Nos. 3,886,057 and 4,050,995 disclose a method for determining water vapor transmission rate by mass spectrometer calculations of oxygen and hydrogen generated during electrolysis. Alternatively, the water vapor transmission rate can be determined by the electrical current needed to electrolyze the water vapor. All of the above-described devices do not account for diffusion resistance and albedo of crops. The need for a device which accounts for these parameters is expressed in Rosenberg, Microclimate: The Biological Environment, John Wiley and Sons, New York, Copyright 1974, page 198 of 315."}
{"text": "Presently, most passenger-side air bag doors are formed in an air bag cover that is separate from the instrument panel because the latter is made of a commercially available thermoplastic material particularly suited to meeting the requirements of its application but is not well suited to meeting the requirements of an air bag door that is defined by a tear seam formed in the instrument panel. For example, the plastics materials used to make a self-supporting instrument panel must have a certain degree of stiffness and high heat resistance to meet the requirements of its application but the materials that are currently available for such an instrument panel do not retain ductility and become embrittled at very low or cold temperatures which is unsuitable for air bag deployment where the air bag door is formed integral with and defined by a tear seam formed in the instrument panel. Styrene-maleic anhydride, polypropylene, polycarbonate and polyphenylene oxide are examples of thermoplastic materials that are suitable for the instrument panel but do not have the required ductility or flexibility for such an air bag door at very low temperatures. As a result, the air bag cover or door on the passenger side is normally made separate from the instrument panel and of a different commercially available thermoplastic material such as polyurethane elastomer, polyester elastomer, and polyolefin elastomer which are suitable for such application but are not suitable for the requirements of the instrument panel.\nCosts can be reduced, quality can be improved and styling can be enhanced by molding the instrument panel and the air bag cover including an integral air bag deployment door in one piece at the same time out of the same commercially available material provided, however, that the normal material requirements for the instrument panel are not sacrificed and the safe assured operability of the integral air bag deployment door is retained at reasonable cost, even though the plastics material may not be well suited thereto."}
{"text": "In recent years, a multiband antenna that enables communication over a plurality of frequency bands for ensuring communication capacity has been put to practical use, for example, as a base station for mobile communication and an antenna device for Wi-Fi (registered trademark) communication equipment.\nA multiband antenna is disclosed as a multiband antenna array in, for example, FIGS. 11, 12, and 13 of PTL 1. This antenna array includes an antenna reflector, and an array of high-band and low-band crossed-dipole antenna elements alternately aligned on the antenna reflector. Further, central conductive fences 1340 are provided between the alignments to reduce mutual coupling. In addition, in PTL 2, in FIGS. 2a, 2b, and 2c, a decoupling element (decoupling structural element 17) is arranged between individual antenna elements (radiation element module 1) of a single-frequency antenna array."}
{"text": "1. Object of the Invention\nThis invention relates to a method for primary cementing a well using a composition of matter containing a polymeric compound which is useful as a drilling mud and may be converted to a cement upon irradiation with a proper source.\n2. Field of the Invention\nThe process of searching for oil and gas is fraught with risk. Approximately three out of every four wells drilled in the United States are dry holes. Even in the instance when a well is found to have penetrated a subterranean formation capable of producing an economic amount of hydrocarbon, the well must be carefully completed after drilling has ended or less than the maximum amount of hydrocarbon will be produced. One problem caused by the improper well completion step of cementing is subterranean movement of gas from a high pressure formation to another formation of lower pressure. The gas lost in this way may never be recovered. This invention solves many of the problems associated with poor cementing procedures by converting the fluid known as \"drilling mud\" directly into a hardened cement. Drilling mud is the fluid typically used during the drilling of a well to lubricate and cool the bit as well as remove rock cuttings from the borehole. Drilling mud is usually displaced in a discrete step by a cement slurry after the borehole is lined with steel casing.\nThe process of drilling a well followed by the steps of casing and cementing it are described below.\na. Drilling the Well\nIn conventional rotary drilling, a borehole is advanced down from the surface of the earth (or bottom of the sea) by rotating a drill string having a drill bit at its lower end. Sections of hollow drill pipe, usually about 30 feet long, are added to the top of the drill string, one at a time, as the borehole is advanced in increments.\nIn its path downward, the drill bit may pass through a number of strata before the well reaches the desired depth. Each of these subsurface strata has associated with it physical parameters, e.g., fluid content, hardness, porosity, pressure, inclination, etc., which make the drilling process a constant challenge. Drilling through a stratum produces significant amounts of rubble and frictional heat; each of which must be removed if efficient drilling is to be maintained. In typical rotary drilling operations, heat and rock chips are removed by the use of a liquid known as drilling mud. As noted above, drill strings are usually made up of sections of hollow pipe terminated by a drill bit. Drilling mud is circulated down through the drill string, out through orifices in the drill bit where the mud picks up rock chips and heat, and returns up the annular space between the drill string and the borehole wall to the surface. The mud is sieved on the surface, reconstituted, and pumped back down the drill string.\nDrilling mud may be as simple in composition as clear water, but more likely it is a complicated mixture of clays, thickeners, and weighting agents. The characteristics of the drilled geologic strata and, to some extent, the nature of the drilling apparatus determine the physical parameters of the drilling fluid. For instance, the drilling mud must be capable of carrying the rock chips to the surface from the drilling site. Shale-like rocks often produce chips which are flat. Sandstones are not quite so likely to produce a flat chip. The drilling fluid must be capable of removing either type of chip. Conversely, the mud must have a viscosity which will permit it to be circulated at high rates without excessive mud pump pressures.\nIn the instance where a high pressure layer, e.g., a gas formation, is penetrated, the density of the drilling mud must be increased to the point such that the hydrostatic or hydraulic head of the mud is greater than the downhole (or \"formation\") pressure. This prevents gas leakage out into the annular space surrounding the drill pipe and lowers the chances for the phenomenon known as \"blowout\" in which the drilling mud is blown from the well by the formation gas. Finely ground barite (barium sulfate) is the additive most widely used to increase the specific gravity of drilling mud; although, in special circumstances, iron ore, lead sulfide ferrous oxide, or titanium dioxide may also be added.\nIn strata which are very porous or are naturally fractured and which have formation pressures comparatively lower than the local pressure of the drilling mud, another problem occurs. The drilling fluid, because of its higher hydrostatic head, will migrate out into the porous layer rather than completing its circuit to the surface. This phenomenon is known as \"lost circulation\". A common solution to this problem is to add a lost circulation additive such as gilsonite.\nFluid loss control additives may be included such as one containing either bentonite clay (which in turn contains sodium montmorillonite) or attapulgite, commonly known as salt gel. If these clays are added to the drilling mud in a proper manner, they will circulate down through the drill string, out the drill bit nozzles, and to the site on the borehole wall where liquid from the mud is migrating into the porous formation. Once there, the clays, which are microscopically plate-like in form, form a filter cake on the borehole wall. Polymeric fluid control agents are also well known. As long as the filter cake is intact, very little liquid will be lost into the formation.\nThe properties required in drilling mud constantly vary as the borehole progresses downward into the earth. In addition to the various materials already noted, such substances as tannin-containing compounds (to decrease the mud's viscosity), walnut shells (to increase the lubricity of the mud between the drillstring and the borehole wall), colloidal dispersions, e.g., starch, gums, carboxy-methyl-cellulose (to decrease the tendency of the mud to form excessively thick filter cakes on the wall o: the borehole), and caustic soda (to adjust the pH of the mud) are added as the need arises.\nThe fluid used as drilling mud is a complicated mixture tailored to do a number of highly specific jobs.\nOnce the hole is drilled to the desired depth, the well must be prepared for production. The drill string is removed from the borehole and the process of casing and cementing begins.\nb. Casing and Cementing the Well\nIt should be apparent that a well that is several thousand feet long may pass though several different hydrocarbon producing formations as well as a number of water producing formations. The borehole may penetrate sandy or other unstable strata. It is important that in the completion of a well each producing formation be isolated from each of the others as well as from fresh water formations and the surface. Proper completion of the well should stabilize the borehole for a long time. Zonal isolation and borehole stabilization are also necessary in other types of wells, e.g., storage wells, injection wells, geothermal wells, and water wells. This is typically done, no matter what the type of well, by installing metallic tubulars in the wellbore. These tubulars known as \"casing\", are often joined by threaded connections and cemented in place.\nThe process for cementing the casing in the wellbore is known as \"primary cementing\". In an oil or gas well, installation of casing begins after the drill string is \"tripped\" out of the well. The wellbore will still be filled with drilling mud. Assembly of the casing is begun by inserting a single piece of casing into the borehole until only a few feet remain above the surface. Another piece of casing is screwed onto the piece projecting from the hole and the resulting assembly is lowered into the hole until only a few feet remain above the surface. The process is repeated until the well is sufficiently filled with casing.\nA movable plug, often having compliant wipers on its exterior, is then inserted into the top of the casing and a cement slurry is pumped into the casing behind the plug. The starting point for a number of well cements used in that slurry is the very same composition first patented by Joseph Aspdin, a builder from Leeds, England, in 1824. That cement, commonly called Portland cement is generally made up of:\n______________________________________ 50% Tricalcium silicate 25% Dicalcium silicate 10% Tricalcium aluminate 10% Tetracalcium aluminoferrite 5% Other oxides. ______________________________________\nAPI Class A, B, C, G and H cements are all examples of Portland cements used in well applications. Neat cement slurries may be used in certain circumstances; however, if special physical parameters are required, a number of additions may be included in the slurry.\nAs more cement is pumped in, the drilling fluid is displaced up the annular space between the casing and the borehole wall and out at the surface. When the movable plug reaches a point at or near the bottom of the casing, it is then ruptured and cement pumped through the plug and into the space between the casing and the borehole wall. Additional cement slurry is pumped into the casing with the intent that it displace the drilling mud in the annular space. When the cement cures, each producing formation should be permanently isolated thereby preventing fluid communication from one formation to another. The cemented casing may then be selectively perforated to produce fluids from particular strata.\nHowever, the displacement of mud by the cement slurry from the annular space is rarely complete. This is true for a number of reasons. The first may be intuitively apparent. The borehole wall is not smooth but instead has many crevices and notches. Drilling mud will remain in those indentations as the cement slurry passes by. Furthermore, as noted above, clays may be added to the drilling mud to form filter cakes on porous formations. The fact that a cement slurry flows by the filter cake does not assure that the filter cake will be displaced by the slurry. The differential pressure existing between the borehole fluid and the formation will tend to keep the cake in place. Finally, because of the compositions of both the drilling mud and the cement slurry, the existence of non-Newtonian flow is to be expected. The drilling mud may additionally possess thixotropic properties, i.e., its gel strength increases when allowed to stand quietly and the gel strength then decreases when agitated. The combination of these effects creates boundary layer conditions which hinder the complete displacement of the drilling mud from the annular space.\nSeveral remedial and preventative steps may be taken to assist in removal of drilling mud from the annular space. Long wire bristles or \"scratchers\" may be placed at intervals along the casing string as it is inserted into the hole. These devices have the direct beneficial effect of removing filter cakes from the borehole wall and providing an improved bonding site for the set cement. However, because of the flow characteristics of the cement slurry and the drilling mud, the scratchers are not completely effective in causing the mud to displace.\nSimilarly, a device known as a \"centralizer\" can be added to the casing string to centralize the casing and improve the flow around the string. Although centralizers are helpful in preventing quiescent areas, the borehole does not have a perfect interior surface and dead spots will occur in, e.g., dog-legs in the hole.\nOther methods of aiding in the displacement operation, have been attempted and each has its own benefits and detriments. These methods include preflushes, spacers, additives to reduce drilling mud viscosity, abrasive materials to erode the filter cake, and high apparent viscosity cement to displace drilling mud in a piston-like motion. None of the known methods is completely effective in removing mud from the borehole wall.\nOne goal of the art has been to dispense with the necessity of displacing the drilling mud by utilizing but a single fluid capable of performing the functions of both the mud and the cement. The benefits of such a multifunctional fluid are apparent. The requirements that the filter cake be removed from the borehole wall and that the mud be taken from the imperfections in the wall are therefore obviated. A few patents disclose methods for converting drilling mud to cement and are discussed below. However, none of these disclosures suggest a drilling mud which is used per se as the well cement. The disclosures normally teach the addition of some other material to the mud prior to its use as a cement. None of the references shows the conversion of mud to cement by irradiation in the well.\nA method of converting drilling mud to a cement slurry is disclosed in U.S. Pat. No. 3,168,139, issued to Kennedy et al on Feb. 2, 1965. Kennedy et al teaches the straightforward step of adding a hydraulic cement, preferably a Portland cement, to the drilling mud to form a cement slurry. Kennedy et al does not suggest using the cement slurry so-formed as a drilling fluid. Kennedy et al also notes at Column 13, line 32 et seq that \"less channeling through the set cement occurs than when conventional cementing slurries are employed\". The efficiency of the Kennedy et al process is therefore not assured.\nThe patent to Cunningham et al, U.S. Pat. No. 3,409,093, issued Nov. 5, 1968, teaches the use of a known cement slurry as the drilling fluid. This process is said simultaneously to produce an impenetrable filter cake on the borehole wall and a strong cement sheath on the walls capable of stabilizing the borehole wall. It should be apparent that close control of setting rate by retarder concentration and water content (via inclusion of water and water-loss additives) must be maintained using this process lest the cement slurry set and seize the drill string. It should also be apparent that the drilling fluid disclosed by Cunningham et al will be useful as a drilling fluid only for a short time.\nTragesser, U.S. Pat. No. 3,557,876, issued Jan. 26, 1971, teaches a drilling mud which can be converted, when desired, to a cementitious material by the addition of an alkaline earth oxide such as calcium, strontium or barium oxide. The particular drilling mud involved comprises water, collodial clay, various conventional additives, and a substance known as pozzolan. Pozzolan is a siliceous material (generally about 50 percent silicon oxide) containing various percentages of other oxides such as magnesium oxide, aluminum oxide, or iron oxide. These materials are said to form a cementitious material when reacted with an alkaline earth oxide in the presence of water at the temperatures found downhole in a well. There is no assurance that adequate mixing between the disclosed drilling mud and the alkaline earth oxide will occur in the well. Without proper mixing, quiescent volumes may remain within the well and prevent attainment of an acceptable cement job.\nThe disclosure in Harrison et al, U.S. Pat. No. 3,605,898 issued Sept. 20, 1971, relates to a hydraulic cement composition containing a setting retardant known as a heptolactone, preferably D-gluco-D-guloheptolactone. The composition is said to be usable as a drilling mud as long as the retardent is effective. A water soluble polyvalent metal salt, preferably CaCl.sub.2, is added to the mud to effect a conversion into cement.\nThe teachings in Miller et al, U.S. Pat. No. 3,887,009, issued June 3, 1975, relate to a clay-free magnesium-salt drilling fluid. Such fluids are said to typically contain 15 to 60 pounds of magnesium sulfate per barrel of drilling mud, from 20 to 70 pounds per barrel dolomite, about 3 to 15 pounds of calcium oxide per barrel, and 4 to 10 pounds of gypsum per barrel. In order to form a cement from this drilling mud, sufficient magnesium oxide, magnesium sulfate and dolomite or magnesium carbonate are added to produce a magnesium oxysulfate cement. Again, control of the concentration of the added material appears to be critical.\nNone of the disclosures of Tragesser, Harrison et al or Miller et al eliminate the step of displacing the drilling fluid from the well.\nThe process disclosed herein does not require the sequential addition of various materials as the time for cementing approaches, although it is acceptable to do so. The step which initiates the conversion of the disclosed composition from drilling mud to cement is the irradiation of that composition in the well. The cross-linkable polymers contained in the polymeric composition thereafter crosslink forming a strong set cement.\nThere are, of course, other known compositions containing polymerizable compounds which have been placed in wells for a variety of reasons.\nA disclosure of one such composition is found in Perry et al, U.S. Pat. No. 3,114,419 issued Dec. 17, 1963. Perry et al suggests a \"method for polymerizing liquid-resin forming materials suitable for use in well bores penetrating permeable subterranean formations\". The preferred method uses radiation to copolymerize an alkylidene bisacrylamide with an ethylenic monomer. Perry et al teaches that the polymeric composition is made up so as to have a specific gravity between 1.07 and 1.18 (Column 4, lines 48 et seq). The specific gravity may be adjusted by the addition of non-ionizing organic weighting agents such as sugar or glycerol. The composition should retain some water solubility to be effective (Column 3, lines 43 et seq). The composition is pumped to the formation to be plugged by first pumping fresh water down the casing followed by the polymeric composition (Column 5, line 65 et seq). An amount of salt water is then pumped in until the composition is placed at the site of the porous formation. Because the specific gravity of the composition is between that of the overlying fresh water and the underlying salt water, it will stay in place. A radioactive source is then inserted into the well, to effect copolymerization and seal the permeable formation. Perry et al additionally discloses the use of similar processes to seal permeable formations when using a gas as the drilling fluid (Column 6, line 46 et seq and Example V) and to plug a permeable formation in a water flood injection well (Example VI).\nPerry et al does not disclose or suggest the use of drilling muds containing polymeric materials which can be converted to cement by irradiation in the well.\nOther disclosures relating to the use of irradiated polymers in well liquids can be found in U.S. Pat. Nos. 3,830,298, 3,872,923, 3,877,522, and 3,973,629 each issued to Knight et al. Another disclosure of radiation induced polymerization may be found in Canadian Pat. No. 1,063,336 to Ressaine et al. Each of these disclosures relates to a particular use of a polymerized \"acrylamide and/or methacrylamide and acrylic acid, methacrylic acid, and/or alkali metal salts thereof\" in a number of different ways, e.g., as a water loss additive in cement, as a plugging medium in a porous formation, etc. However, each of the disclosures deals with a polymeric composition which is irradiated prior to being placed in a well.\nc. Completing the Well\nAfter the casing is cemented in place, a well is prepared for hydrocarbon production using a number of separate steps. A perforation tool is commonly lowered within the cemented casing to the region of a producing formation. The perforation tool is a device which often is constructed of a number of guns which produce holes through the casing and its enclosing cement and into the producing formation. The interior of the casing is thereafter in open communication with the formation.\nOther completion steps may involve fracturing to increase well productivity, installing screens to exclude sand from the well bore, and installing production tubing between the producing formation and the surface."}
{"text": "In contemporary vehicular alternator system voltage regulators, it is commonplace to use temperature dependent voltage sources to synthesize battery charge profiles. These temperature dependent charge profiles are necessary, because these systems are used to charge batteries that can be most efficiently charged at different rates depending on ambient temperature proximate the battery. For instance, some of these systems rely on generation of reference signals that have substantially different temperature coefficients that are used to determine thresholds used in regulation and fault detection operations.\nPrior art systems accomplish this requirement by independently synthesizing two temperature dependent reference signals, or voltages. These voltages are used independently or in combination to synthesize the appropriate charging and fault detection profiles, and associated thresholds, for the battery under charge. In FIG. 1 a typical prior art block diagram is illustrated. A temperature dependent voltage source 101 has an output 103 for providing a voltage having a temperature dependence of -680 ppm/.degree.K. (parts per million/degree Kelvin) which is equivalent to -10 mV/.degree.K., (millivolts/degree Kelvin) which is typical in a 12 volt battery charging system. Also, a second voltage source 105 provides an output 107 for generating a temperature independent voltage--here 0 ppm/.degree.K., or 0 mV/.degree.K. The outputs 103 and 107 are supplied to a charge regulator 109 which in turn controls an alternator 104 to generate a temperature dependent voltage to a battery 113. The temperature dependent voltage provides a charge profile dependent on the voltage provided at the output 103. Also the charge regulator 109 has a fault detection function dependent on the voltage provided at the output 107.\nTypically, two bandgap voltage sources are used to construct these temperature predictable voltage sources. This is the cause of several problems. First, since there are two complete bandgap reference circuits, the size; complexity; and cost of the combined circuits is unnecessarily high. Secondarily, because these voltage sources 101 and 105 operate independently, the charge regulator 109 cannot expect the outputs 103 and 107 of the two voltage sources 101 and 105 to be accurately related over temperature. This is problematic because the charge regulator 109 requires that the outputs 103 and 107 of the two voltage sources 101 and 105 be closely related, in order to minimize the error of the desired charge profile and associated thresholds. If the outputs 103 and 107 of the two voltage sources 101 and 105 are not correlated, then the charge regulator 109 must be at least as inaccurate as the independence between the two voltage sources 101 and 105.\nFurthermore, to improve any correlation between the two voltage sources 101 and 105, active trimming of critical components is required. This is process intensive and expensive to perform in a manufacturing environment. At a minimum, this dual active trimming process adds unnecessary circuit area and test time.\nWhat is needed is an improved dual output temperature compensated voltage reference, that has a known voltage relationship over temperature between its dual outputs, is simpler, more accurate, more compact in size, lower in cost, and is more manufacturable."}
{"text": "Field of the Invention\nThe present invention relates to a method of and an apparatus for interengaging a pair of slide fastener stringers equipped with top stops, sliders, separable bottom stops before the stringers are cut to individual slide fastener lengths."}
{"text": "In recent years, much attention has been paid to problems of environmental destruction on global scale such as global warming, ozone depletion, acid rain, etc. and it is urged to internationally take countermeasures against global pollution. Along with this problem, various regulations have been set for environmental conservation. Under such circumstances, it is a serious problem to discharge organic solvents (VOC) into the atmosphere, and in each industrial field, along the trend for tightening regulations against VOC, it is promoted to seek alternatives to organic solvents (alternatives to VOC). In the paint industry, as an alternative to a conventional organic coating material, a powder coating material is expected to be promising as a coating material which contains no VOC and which is environmentally friendly such that it does not require exhaust treatment or waste water treatment, and it can be recovered for reuse.\nAs a powder coating material, an acrylic resin type, a polyester resin type and an epoxy resin type have been mainly used. However, there is a problem that such coating materials are poor in weather resistance.\nOn the other hand, as a coating material to solve such a problem, a fluorinated resin powder coating material comprising a fluorinated resin has been developed (for example, Patent Document 1).\nFurther, for reducing the cost, etc. a powder coating material in which a polyester resin and a fluorinated resin are blended has been proposed (for example, Patent Document 2).\nAs compared to the acrylic resin type, polyester resin type and epoxy resin type powder coating materials, the powder coating material in which a polyester resin and a fluorinated resin are blended is excellent in weather resistance, however, depending on a kind of a polyester to be used, it is difficult to obtain the effect to improve the weather resistance derived from a fluorinated resin layer in some cases.\nFurther, by coating a polyester resin with a fluorinated resin to form a double coating, a coating film which is excellent in weather resistance can be formed at a low cost. However, steps for forming the coating film are complicated since it is necessary to carry out double coating."}
{"text": "Decoys are widely manufactured and sold for the purpose of attracting wild animals for gaming and observation. “Gaming” includes such activities as hunting, and “observation” includes such activities as photography.\nIt has been empirically observed that some game animals are more attracted to decoys having movement than they are to decoys that remain still. For this reason, there is a need for apparatus that provide motion to conventional decoys in a manner that helps to attract wild animals."}
{"text": "In order for the human eye to have clear vision of objects at different distances, the effective focal length of the eye must be adjusted to keep the image of the object focused as sharply as possible on the retina. This change in effective focal length is known as accommodation and is accomplished in the eye by varying the shape of the crystalline lens. Generally, in the unaccommodated emmetropic eye the curvature of the lens is such that distant objects are sharply imaged on the retina. In the unaccommodated eye near objects are not focused sharply on the retina because their images lie behind the retinal surface. In order to visualize a near object clearly, the curvature of the crystalline lens is increased, thereby increasing its refractive power and causing the image of the near object to fall on the retina.\nThe change in shape of the crystalline lens is accomplished by the action of certain muscles and structures within the eyeball or globe of the eye. The lens is located in the forward part of the eye, immediately behind the pupil. It has the shape of a classical biconvex optical lens, i.e., it has a generally circular cross section having two convex refracting surfaces, and is located generally on the optical axis of the eye, i.e., a straight line drawn from the center of the cornea to the macula in the retina at the posterior portion of the globe. In the unaccommodated human eye the curvature of the posterior surface of the lens, i.e., the surface adjacent to the vitreous body, is somewhat greater than that of the anterior surface. The lens is closely surrounded by a membranous capsule that serves as an intermediate structure in the support and actuation of the lens. The lens and its capsule are suspended on the optical axis behind the pupil by a circular assembly of very many radially directed elastic fibers, the zonules, which are attached at their inner ends to the lens capsule and at their outer ends to the ciliary muscle, a muscular ring of tissue, located just within the outer supporting structure of the eye, the sclera. The ciliary muscle is relaxed in the unaccommodated eye and therefore assumes its largest diameter. According to the classical theory of accommodation, originating with Helmholtz, the relatively large diameter of the ciliary muscle in this condition causes a tension on the zonules which in turn pulls radially outward on the lens capsule, causing the equatorial diameter of the lens to increase slightly and decreasing the anterior-posterior dimension of the lens at the optical axis. Thus, the tension on the lens capsule causes the lens to assume a flattened state wherein the curvature of the anterior surface, and to some extent the posterior surface, is less than it would be in the absence of the tension. In this state the refractive power of the lens is relatively low and the eye is focused for clear vision for distant objects.\nWhen the eye is intended to be focused on a near object, the ciliary muscles contract. According to the classical theory, this contraction causes the ciliary muscle to move forward and inward, thereby relaxing the outward pull of the zonules on the equator of the lens capsule. This reduced zonular tension allows the elastic capsule of the lens to contract causing an increase in the antero-posterior diameter of the lens (i.e., the lens becomes more spherical) resulting in an increase in the optical power of the lens. Because of topographical differences in the thickness of the lens capsule, the central anterior radius of curvature decreases more than the central posterior radius of curvature. This is the accommodated condition of the eye wherein the image of near objects falls sharply on the retina.\nPresbyopia is the universal decrease in the amplitude of accommodation that is typically observed in individuals over 40 years of age. In the person having normal vision, i.e., having emmetropic eyes, the ability to focus on near objects is gradually lost, and the individual comes to need glasses for tasks requiring near vision, such as reading.\nAccording to the conventional view the amplitude of accommodation of the aging eye is decreased because of the loss of elasticity of the lens capsule and/or sclerosis of the lens with age. Consequently, even though the radial tension on the zonules is relaxed by contraction of the ciliary muscles, the lens does not assume a greater curvature. According to the conventional view, it is not possible by any treatment to restore the accommodative power to the presbyopic eye. The loss of elasticity of the lens and capsule is seen as irreversible, and the only solution to the problems presented by presbyopia is to use corrective lenses for close work, or bifocal lenses, if corrective lenses are also required for distant vision.\nCertain rings and/or segments have been used in ocular surgery for various purposes. Rings and/or segments of flexible and/or elastic material, attached or prepared in situ by fastening the ends of strips of the material around the posterior portion of the globe, posterior to the pars plana (over the underlying retina), have been used to compress the sclera in certain posterior regions. Supporting rings of metal, adapted to fit the contour of the sclera have been used as temporary supporting structures during surgery on the globe. However, none of these known devices have been used for surgical treatment of presbyopia, and none have been adapted to the special needs of prosthetic devices used in treating presbyopia.\nAccordingly, a need has continued to exist for a method of treating presbyopia that will increase the amplitude of accommodation of the presbyopic eye, thereby lessening or eliminating the need for auxiliary spectacle lenses to relieve the problems of presbyopia."}
{"text": "This invention relates to thermal printing and more particularly to means for assembling and making connections to a thermal printhead.\nThermal printing presents many advantages over mechanical means, particularly when implemented with electronic calculators. These advantages include lower cost of equipment and fewer moving parts, resulting in fewer mechanical failures. When a printed tape showing keyboard entries and calculated results is desired, a thermal printing system may be utilized to produce such an output. Using thermally sensitive paper which may be advanced as in conventional printers, a thermal printing system may accept signals from the calculator, decode those signals, use the decoded signals to selectively heat portions of the printhead, and thereby produce characters or numerals on heat sensitive paper.\nWith reference to FIG. 1, thermal printheads are generally produced with a plurality of heating elements or mesas 16 which may be selectively energized to produce a pattern of lines or dots. When the proper signals are provided to these elements or mesas 16 by means of metallized leads combined with proper advancement of the thermally sensitive paper 17 by means of a platen-advancing mechanism 18, characters or numerals are formed on the thermally sensitive paper 17. The heating elements 16 must be in very close contact with the paper 17 for satisfactory printing. These elements 16, as part of a thermal printhead assembly 14, are typically held in contact with the thermally sensitive paper 17 by a spring-loaded pivot arrangement 15.\nThe heating elements or mesas 16 are located on the one side of a substrate 10 which is usually ceramic or a similar material. The leads are generally brought through the ceramic substrate to the other side where pins 11 are attached. This combination of the ceramic substrate 10, the heating elements or mesas 16, and the leads 15 hereinafter described generally as a printer ceramic.\nA standard lead frame 12 may be provided in which to insert the pins 11. From the lead frame 12, individual wire conductors 20 may be used to make connections to the printer logic 21. Because there is usually very little clearance between the substrate 10 and the platen 18 which advances the thermally sensitive paper 17 (not enough to locate a standard lead frame or other conventional connector), the leads from the heating elements 16 must be brought through the ceramic substrate 10. The lead frame 12 not only allows connection to the heating elements or mesas 16 but may also operate to hold the printer ceramic in the proper position. The printhead may also be held in place by attaching the ceramic substrate 10 to a pivoting bracket 15. A number of printer ceramic-lead frame assemblies may be installed adjacent to each other, in a printhead assembly 14, thereby providing for the printing of as many digits or characters as are desired.\nThere are a number of difficulties with such an arrangement. The process of bringing metal leads from the heating elements 16 through the ceramic substrate 10 and attaching pins 11 to them is more expensive than merely providing metallized leads on a single side of the substrate 10. In addition, the lead frames 12 are bulky and require a good amount of space in an area where space is at a premium. Further, the use of printer ceramic-lead frame assemblies often makes replacement of printer ceramics a difficult task since a housing 13 may first have to be disassembled. The attachment of the ceramic substrate 10 to a pivoting bracket 15 may require removal of screws, bolts, or other fasteners 19. Finally, because so much space on the ceramic substrate 10 is required by the pins 11 and lead frame 12, there is little room for thermal contact between the substrate 10 and metal parts of the thermal printer system which might act as heat sinks. The thermal contact that is available is some distance away from the heating elements or mesas 16 themselves, and the ceramic is not a good conductor of heat. The heating elements 16 of the thermal printer chip must cool quickly so that, as the thermally sensitive paper 17 is advanced, there will be no trailing effect. Unless proper heat sinking is available, the elements 16 will still be hot enough when the paper is advanced to leave a streak on the thermally sensitive paper 17. To avoid this problem, printing must be slowed down to allow time for the elements 16 to cool before the paper 17 is advanced, or heat sinking must be provided. Thus, proper heat sinks will allow faster operation."}
{"text": "The present invention relates generally to methods and apparatus for treating fluids, and specifically relates to methods and apparatus for treating fluids through contact with a contact media. The method and apparatus of the present invention are particularly suitable for use in ion exchange operations where a fluid is contacted with a fluidizable ion exchange media such as a resin, or with filtration or absorption operations where the fluids contacted with a media such as activated carbon.\nIon exchange processes for testing fluids are well known. Such ion exchange operation may include, for example; \"water softening\", deionization, de-alkylizing, desilicizing, and organic scavenging. With respect to the present invention, these ion exchange processes will for convenience be discussed in terms of water treatment. It should be clearly understood, however, that the methods and apparatus of the present invention may be utilized in the treatments of many other fluids.\nTypically, an ion exchange process is affected by flowing the water through a vertical column of an ion exchange contact media, typically a resin. As the water contacts the resin, ions in the water will be attracted to the resin from the water. One type of resin may be utilized to remove cations from the water and a second type of resin may be utilized to remove anions from the water. Preferably, the separate resins will be contained in separate beds. However, conventional techniques of water deionization include the use of two resins mixed in a single bed.\nThe softening of water by ion exchange is accomplished by replacing the calcium and magnesium ions in the water by an equivalent number of sodium ions from the resin. Resin in the bed will contain only a finite number of exchangeable sodium ions. This number defines the \"capacity of the resin\". When the capacity of the resin has been exhausted, i.e., when all of the exchangeable sodium ions on the resin have been replaced by calcium and magnesium ions from the water, the resin must be regenerated back to the sodium form. This regeneration is typically accomplished by passing a sodium chloride solution (a brine) through the resin. Additionally, the resin will be rinsed to remove excess brine, and will be backwashed to remove particulate matter which may have accumulated in the resin during the ion exchange step (\"the service cycle\").\nAs water flows through a bed of resin, the majority of the ion exchange will take place in the portion of the resin which is first contacted by the fluid. In an ion exchange system where the fluid flows downwardly through the bed, this exchange creates an \"exhausted band\" of exhausted resin which expands downwardly through the bed as the operation continues. When the band approaches the bottom of the resin bed, the bed must be regenerated as discussed above.\nAdditionally, a vertical column of resin operating in a service cycle has an exchange zone, or active band, starting at the top and moving down through the bed of resin. The width of the active band varies with certain operating parameters of the system. For example, as the service flow is increased the active band will spread out. The resin bed must be removed from service and regenerated before the active band reaches the bottom of the column to prevent leakage of the ions being removed from the fluid. This prevents full utilization of the resin since there is resin not fully exhausted in and below the active band. When a column of resin in normal operation is regenerated, sufficient regenerant (salt in the case of water softening) is used to regenerate the entire volume regardless of what percent of the bed was actually exhasted.\nMany conventional vertical ion exchange columns are designed to function both as a column for the service cycle, i.e., for the initial ion exchange process, and for the regeneration cycle, and therefore also include provisions for backwashing, regeneration, and rinsing of the resin. This structure requires the fluid influent to be shut off from the column while the regeneration operation takes place, thereby interrupting the supply to service of treated fluid. When an uninterrupted supply of treated fluid is required, a second vertical column is typically provided. This second column will be regenerating during the service cycle of the first column and vice versa. Conventional columns typically include a shell to contain the resin, a support for the resin, and means for distributing flow both upwardly and downwardly through the resin, (for both the service cycle and the regeneration and backwashing cycles). The shell must have sufficient space above the resin to allow the resin to expand during the backwashing operation. Valves and controls are typically necessary to bypass \"raw\" (untreated) fluid around the column during the regeneration cycle, to inject the regeneration fluid into the column, and to reverse the direction of fluid flow for backwashing.\nConventional vertical columns may include several disadvantages. Where a single large resin bed is utilized, fluid has a tendency to channel through the bed during periods of low flow rate, thereby reducing the effective contact of the fluid with the resin. Additionally, the requirement of additional space above the resin bed to facilitate the backwashing operation adds cost to the vessel. Where uninterrupted service is required and a second unit is provided, the additional unit adds significant cost and size to the unit. A significant factor in this cost is that a control valve must be provided to switch fluid flow from one vessel to the other. This control valve must be large enough to provide a significant flow of the influent into the column without placing an excessive pressure drop in the system. Large control valves of this type typically contribute a significant portion of the cost of conventional ion exchange units. These valves still often place an undesireable pressure drop in the system.\nIf only one column is provided, in typical conventional systems, not only must the flow of treated water be interrupted, but untreated water must be used with the unit itself for the backwashing, regeneration, and rinsing operations. This use of untreated water will, in itself, decrease the operating efficiency of the contact media regeneration, and will therefore similarly decrease the efficiency of the ion exchange process.\nBecause of the deficiencies discussed above, several attempts have been made to devise methods and apparatus for an uninterrupted or, continuous ion exchange process in a single column. Typically, these processes involve the movement of the contact media downwardly through the column or ion exchange vessel while the fluid flows upwardly through the column. In some cases, the resins are actually fluidized, or suspended, in the fluid flow. This upward flow, and especially fluidization, typically provide less than optimal ion exchange. A major factor in the efficiency of ion exchange process is the physical contact of the water molecules with the resin. With a downward fluid flow, both gravity and the influent flow serve to compact the resin into a tightly formed bed. This compacting of the bed forces the fluid to flow closer to the resin beads, causing surface effects on the water and forcing fluid to flow into the pores of the beads. This compacting of the bed, therefore, increases both contact efficiency and the bed capacity. In contrast, an upward flow, as found in the prior art, causes the resin to expand, as noted above, sometimes to the point of fluidization. This unpacked state of the resin causes a reduction in contact with the water. Additionally, even when a system is designed to operate with the resin not in a state of fluidization, flow rate surges must be prevented to avoid the fluidization.\nAs indicated above, conventionally proposed continuous ion exchange methods and apparatus typically move the resins downwardly through the exchange vessel. When each portion of the resin reaches a predetermined location in the vessel, the resin is removed and regenerated in a separate vessel. Therefore, for efficient use of the system and the contact media, the rate of travel of the media must be regulated in response to the rate of flow of the influent.\nConventionally proposed methods and apparatus for continuous ion exchange typically provide for contacting the fluid with a single resin in a vessel. However, for operations such as water deionization or demineralization, the water is typically contacted with two ion exchange resins in two stages. In the first stage, the water contacts a first resin which will attract the cations from the water and replace them with hydrogen ions. This first resin is typically regenerated with an acid. In the second stage, the fluid is contacted with a second resin which attracts the anions from the water and replaces them with hydroxide ions. This second resin is typically regenerated with a base, such as sodium hydroxide. The hydrogen ions from the first stage and the hydroxide ions from the second stage combine to produce water. Similarly, fluid may be contacted with selected resins for other treatment operations, such as dealkylization, etc. Each of these resins will require a selected regenerization fluid. The fluid must, therefore, be treated with the different resins in different columns or with a mixture of the two resins. If a mixture of the resins is utilized, then the resins must be separated prior to regeneration.\nAdditionally, activated carbon is often utilized in fluid treatment operation. The activated carbon may be used to remove gases and other organic impurities giving taste and odor to drinking water. Activated carbon may be used as a pretreatment for water supplied to water softening systems to prevent organic fouling of the resin beds, and may also be utilized in waste water treatment. Although activated carbon is not susceptible to being \"regenerated\" per se, the carbon periodically requires backwashing to remove suspended matter and to re-grade the bed. With conventionally proposed continuous water treatment operations, an activated carbon section must be contained within a separate vessel.\nAccordingly, the present invention provides a new method and apparatus for the continuous treatment of fluids in a single column wherein the contact media need not be treated until it is completely exhausted. Additionally, multiple contacting media may be utilized for different types of ionic exchange, absorption or filtration in the column, and each may be regenerated without interrupting the continuous service flow. Also, this continuous fluid treatment can be performed with an optimal downward fluid flow. The invention, thus, significantly overcomes the deficiencies presented by the prior art."}
{"text": "The present invention relates to a weight measuring device and, more particularly, to a weight measuring device, for a heating appliance such as a microwave oven, which measures a weight of an object by detecting a voltage produced from a piezoelectric transducer element.\nA microwave oven including a weight measuring device has been developed for automatically heating a food based on the weight of the food. For example, a weight measuring device is disclosed in U.S. patent application Ser. No. 724,072 filed on Apr. 16, 1985, entitled \"FOOD WEIGHT MEASURING DEVICE FOR COOKING APPLIANCE\" by Taisuke MORINO et al. The British counterpart is application No. 8509781 filed on Apr. 17, 1985, and the counterpart in West Germany is application No. P3514505.6 filed on Apr. 20, 1985. In Ser. No. 724,072, the weight of the food is measured based on a change in an electrostatic capacity. However, a number of components for the weight measuring device may be for the assembly, thus requiring much labor for assembly, and making it difficult to reduce the costs. The components should be accurately assembled.\nAccordingly, it is desired that a novel weight measuring device be provided to exactly measure the weight of the object based on the output of a piezoelectric transducer element, thereby reducing the number of components."}
{"text": "1. Technical Field\nThis invention relates generally to apparatus and methods for removal of controlled amounts of flowable materials from within sealed containers, and especially to those devices wherein food materials having wide ranges of viscosities, temperatures, and particulate contents are press pumped from flexible walled bags, which are in turn supported within a more rigid container.\n2. Background Art\nFlexible walled bottles have long been used to dispense fluids by means of squeezing. Such bottles have been usually made with resilient walls which tend to restore to their original shape after removal of deforming forces. As the ratio of inside air to liquid increases with bottle use, difficulties are encountered in attempting to completely empty the bottle in a timely manner. Solutions to these problems have included placement of the fluid in an inner bag which is sealed within a squeeze bottle. Examples of such structures have been disclosed by Streck in U.S. Pat. No. 4,865,224 and Uhlig in U.S. Pat. No. 4,098,434.\nAlthough the above solutions decreased the time required to empty the bottle, they led to new problems associated with premature bag collapse with potential exit orifice blockage. Semi-rigid internal bag cartridges which lessen these problems were shown early in U.S. Pat. No. 2,608,320 to Harrison, and various supporting structures placed inside the bag are well known, as recently illustrated by U.S. Pat. No. 5,156,300 to Spahni et al. Combinations of such internal bag elements with dispensing means are available, as typified by U.S. Pat. No. 4,138,036 to Bond wherein a helical coil bag insert is coupled with a dispensing spout supported in the neck opening of a rigid external structure. Other combination devices have mated a protected dispensing orifice with bag cutting and sealing means. An example of this approach is shown by Knorr in U.S. Pat. No. 5,127,550, wherein a first member having a throughbore therein is used in conjunction with a second cutting means to seal adjacent portions of a bag wall together while piercing the bag wall for fluid release through the narrow passages cut therein.\nWhile the above devices have been useful in dispensing relatively nonrepeatable doses of liquids having sufficiently low viscosity to permit fluid flow through given fixed exit orifices, difficulties arise in obtaining repeatable portioning of fluids over wide ranges of temperature and viscosity. In the restaurant industry for instance, such devices have been useful in manual dispensing of condiments by customers.\nIn the area of rapid food preparation however, a need exists to dispense well controlled portions of liquids of various viscosities, or of emulsions, or of colloidal suspensions which can coagulate or settle out with time, or of semi-solid mixtures, such as relish, beans, meat, sour cream, cheese sauces, ketchup, or mustard, all of which have very uneven flow characteristics. In addition, dispensers are required to hold and dispense product at a wide range of temperatures, and must provide maximum safety in dispensing and storage of food products in their original container. All dispensing apparatus used must be approved by the Federal Drug administration (F.D.A.), by the National Sanitation Foundation (N.S.F.) and by local safety regularity authorities. In order to meet such requirements, equipments used must be easily cleaned, made from material that does not support bacterial growth, have minimum bend radii, and eliminate cracks which can hide food particles.\nCurrent trends in the fast food industries indicate continued world wide growth. This will require increased employment of manual labor if the industry continues to hand make the food product sold to their customers. Since this labor force is generally unskilled and has a high turn over rate, the industry will have to automate the food handling both at the restaurant level and at food manufacturing facilities. In addition to labor, cost and space savings, automation will also solve other problems. These include less manual food handling, minimization of incorrect temperature exposure, increased environmental protection, reduced waste, and a more exact ingredient portioning.\nThe restaurant and consumer industries are making increased use of flexible bags to contain food and other products. These bags are easy to handle and require less storage and shipping space. Current semi-automation is achieved by transferring ingredients into mechanical handheld dispensing devices or through the use of infitments. Infitments are attached to the bag and used to interface with a dispense hose or nozzle. These infitments are either placed on the bag before loading or attached at location of use. Both means increase the material cost and require additional labor to attach the infitments or fill the mechanical handheld device. Additionally, the excessive product handling necessary to dispense the ingredients causes undesirable waste and food handling risks.\nExisting technology for automation at the restaurant level consists of equipment to dispense soft drinks, and timing devices for cooking or melting of food products. Some ingredient dispensing devices use peristaltic or diaphragm pumps, but are limited in their ability to accurately change portioning. Pneumatic pumps, using air pressure to discharge product from a bag are currently the only ingredient pumps that meet most of the restaurant automation requirements. For instance, they are not approved for low acid foods or foods considered to be hazardous, such as meat or beans. Pneumatic pump approaches however are limited in their ability to accurately dispense, usually have a valve and or sensors exposed to the food, and require a source of constant or regulated high pressure air. In these devices the portion or amount of ingredient that is pumped is usually dependent on one or more of the following: air pressure, its rise and fall time, duration of pressure application, valves, exit hole size and product viscosity. As the ingredients are displaced from the flexible bag, the volume of air to be pressurized increases, and the time to achieve pressurization becomes longer. If the time that the air pressure is applied is held constant, then the ingredient portion being dispensed will become smaller as the residual ingredient left in the bag is reduced. A need exists therefore, for a dispensing system which will alleviate the above problems by provision of automatic dispensing apparatus which will provide uniform controlled dosages, be adaptable to operate within an overall computer controlled automation system, and meet the demanding cost and safety requirements of the restaurant industry."}
{"text": "1. Field of the Invention\nThe present invention relates to waste treatment systems and especially those waste treatment systems which are packaged or containerized units. Even more particularly, the present invention relates to an air diffuser apparatus in which an air inlet would be provided to convey air into an air container envelope within an aerator body and further allowing air to discharge from the body through tear-shaped openings, and allowing solid particles to be discharged through discharge means on the lower portion of the aerator body.\n2. General Background and Prior Art\nIn the treatment of waste, there is often utilized a containerized or packaged unit treatment plant which treats received waste product on a continuous basis. In the marine industry, for example, there is utilized a package sewage treatment plant which can be used on ships, oil rigs, and like installations. Such sewage treatment devices usually receive flow intermittantly and must treat the flow in order to meet ecological and other marine standards. It is desirable that waste be treated such as sewage in an economical way utilizing as little energy as possible and as few moving parts as possible. Note also it is desirable that a minimum of sludge removal would be required since the sludge disposal presents an extra problem. Further, the accumulation of sludge within the unit would produce an undesirable possibly corrosive attack on the unit itself. An accumulation of sludge within the unit further results in a deterioration of the treatment process.\nIt would be desirable that a sludge treatment apparatus produce a total homogenation of the fluids received from the waste stream in order that they be properly bio-degraded. Waste entering the unit would normally be heterogeneous in nature, containing solid waste material as well as some liquid material. A breakup of this material is necessary in order that it be treated properly.\nIn aerated treatment systems, a problem is faced in that a clogging of the aeration assembly which provides bubbles to the unit will cause a degeneration of the treatment process or in fact a total stoppage of air flow to the vessel, transferring the process from aerobic thus removing any treatment capability. It would be desirable that the aeration unit or air diffuser be of a nature which would reduce or prevent clogging by solid material which may enter the unit.\n3. General Discussion of the Present Invention\nThe present invention would solve prior art problems and shortcomings by providing a waste treatment apparatus having no moving parts in a basic hydraulic system which would consist of an air diffuser assembly having an aerator body with a plurality of side walls connected to form an enclosure and at least a pair of intersecting angularly oriented side walls, forming an air container envelope for retaining an air volume therewithin. There would also be adapted an air inlet which would convey air into said air body at the air container envelope, and a plurality of tear-shaped openings being provided in the body below the air container envelope for allowing air to discharge from the body through tear-shaped openings, with each of the tear-shaped openings providing an upper pointed portion which would minimize the bubble-size of air discharged through the openings. Also, there would be provided a baffle plate mounted within the enclosure at the air inlet, the baffle plate would be used for preventing a short circuit travel of air received by the enclosure to the tear-shaped openings placed on the body adjacent the air inlet. There would also be included an attachment means associated with the air inlet for attaching the aerator body to an air conveying conduit.\nThus, it is an object of the present invention to provide an air diffuser assembly which produces a bubble stream having a high surface area.\nIt is another object of the present invention to provide an aeration diffuser adapted for disallowing the accumulation of sludge or other undesirable deposits within the diffuser assembly, and to disallow the clogging of bubble generating orifices by waste solid material.\nIt is still another object of the present invention to provide a waste treatment apparatus which eliminates the need for sludge removal and further sludge treatment.\nIt is a further object of the present invention to provide a waste treatment system which is gravity operated, utilizing no other energy dependent driving force.\nIt is a further object of the present invention to provide a waste treatment apparatus having distinct liquid operating levels which are maintained without the use of pumps or valves."}
{"text": "The present invention relates to aqueous tertiary di-.beta.-hydroxy amine oxides and particularly to their stabilization against phase separation.\nWieland [Ber., 54, 2353 (1921)] first proposed the oxidation of tertiary amines with hydrogen peroxide. He assumed that the hydrogen peroxide added first to the amine to form an ammonium peroxide compound which then decomposed to amine oxide and water. This intermediate ammonium peroxide complex formation has been substantially by Oswald et al [J. Org. Chem., 28, 657 (1963)]. Trialkylaminehydrogen peroxide adducts have been isolated from reactions using 90+% H.sub.2 O.sub.2 at low temperatures (e.g. -50.degree. C.). These complexes are unstable, colorless liquids at room temperature and have been shown by infrared spectroscopy to be hydrogen bonded polar complexes. These complexes are thermally unstable and decompose to yield trialkylamine oxides and water. This decomposition is quite rapid at 50.degree. C.\nThe most obvious method for the direct oxidation of tertiary amines is to use concentrated hydrogen peroxide sans solvent. However, Hoh et al [J.A.O.C.S., 40, 268 (1963)] has shown that oxidation of dimethyldodecylamine without added solvent at 60.degree. C. using 90% H.sub.2 O.sub.2 proceeds only to 40% completion in 3 hours reaction time. With 70% aqueous H.sub.2 O.sub.2 at 50.degree. C., a 45% completion was realized in 3 hours; however, the completion leveled off at 60% oxide in 10 hours reaction time. Using 35% aqueous hydrogen peroxide, the reaction proceeds more rapidly to yield 50% amine oxide in less than 2 hours reaction time and about 85% amine oxide at the completion of the reaction.\nThe preferred method determined by Hoh et al was the addition of 35% aqueous hydrogen peroxide to the tertiary amine with stirring at 60.degree. C. over a 1 hour period during which time the mixture became gelatinous. Sufficient water was added to keep the reaction mixture fluid. Upon completion of peroxide addition, the proportion of water required to yield a 30%-40% amine oxide solution was added and the temperature was raised to 75.degree. C. A conversion approaching 100% then was realized in 2 hours reaction time. The incremental addition of water was found to give rapid oxidation initially and a quick completion of the reaction. Alternatively, when all of the water was added at the commencement of the reaction, the initial rate of reaction was slow and excessive reaction times were experienced. Hoh et al determined that the optimum reaction temperature for their reaction procedure was about 60.degree.-65.degree. C. They noted that at higher temperatures some decomposition apparently occurred (a yellow color developed) and at lower temperatures the reaction rate decreased. Present day commercial manufacturing of aqueous tertiary amine oxides typically involves the incremental addition of aqueous hydrogen peroxide to the tertiary amine and water.\nAqueous trialkyl, ether dialkyl and polyoxyalkylene amine oxides (e.g. 35%-50% amine oxide solids) appear to provide stable, one phase systems upon storage. However, it has been determined that aqueous tertiary di-.beta.-hydroxy amine oxides are not storage stable but split into two layers upon standing. Many of these hydroxy amine oxides cannot be produced by the noted oxidation reaction without added aqueous solvent and are hygroscopic so that removal of water therefrom is difficult if not impossible. Thus, a great need exists for stabilizing such aqueous amine oxide systems."}
{"text": "The present invention relates to tumor necrosis factor receptor (TNF-R) related polypeptides (proteins) and their ligands, hereinafter referred to as TR1, TR2, TL2 and TL4. The invention also relates to inhibiting or activating the action of such polypeptides using agonists or antagonists by the screening methods described herein.\nMany biological actions are a response to certain stimuli and natural biological processes, and are controlled by factors, such as cytokines. These cytokines act through target cell receptors by engaging the receptor and producing an intracellular response.\nFor example, tumor necrosis factors (TNF) alpha and beta are cytokines which act through TNF receptors to regulate numerous biological processes, including protection against infection and induction of shock and inflammatory disease. The TNF molecules belong to the xe2x80x9cTNF-ligandxe2x80x9d superfamily, and act together with their receptors or counter-ligands, the xe2x80x9cTNF-receptorxe2x80x9d superfamily. So far, ten members of the TNF ligand superfamily have been identified and thirteen members of the TNF-receptor superfamily have been characterized.\nAmong the ligands there are included TNF-a, lymphotoxin-a (LT-a, also known as TNF-b), LT-b (found in complex heterotrimer LT-a2-b), FasL, CD40L, CD27L, CD30L, 4-1BBL, OX40L and TRAIL ((Wiley et al.Immunity 3: 673-682 (1995)) All but one of these (LTa) are expressed as type II membrane proteins. The superfamily of TNF receptors includes the p55TNF receptor, p75TNF receptor, TNF receptor-related protein, FAS antigen or APO-1, CD40, CD27, CD30, 4-1BB, OX40, low affinity p75, NGF-receptor(Meager, A., Biologicals, 22:291-295 (1994)).\nMany members of the TNF-ligand superfamily are expressed by activated T-cells, implying that they are necessary for T-cell interactions with other cell types which underlie cell ontogeny and functions. (Meager, A., supra).\nConsiderable insight into the essential functions of several members of the TNF receptor family has been gained from the identification and creation of mutants that abolish the expression of these proteins. For example, naturally occurring mutations in the FAS antigen and its ligand cause lymphoproliferative disease (Watanabe-Fukunaga, R., et al., Nature 356:314 (1992)), perhaps reflecting a failure of programmed cell death. Mutations of the CD40 ligand cause an X-linked immunodeficiency state characterized by high levels of immunoglubulin M and low levels of immunoglobulin G in plasma, indicating faulty T-cell-dependent B-cell activation (Allen, R. C. et al., Science 259:990 (1993)). Targeted mutations of the low affinity nerve growth factor receptor cause a disorder characterized by faulty sensory innovation of peripheral structures (Lee, K. F. et al, Cell 69:737 (1992)).\nTNF and LT-a are capable of binding to two TNF receptors (the 55- and 75-kd TNF receptors). A large number of biological effects elicited by TNF and LT-a, acting through their receptors, include hemorrhagic necrosis of transplanted tumors, cytotoxicity, a role in endotoxic shock, inflammation, immunoregulation, proliferation and anti-viral responses, as well as protection against the deleterious effects of ionizing radiation. TNF and LT-a are involved in the pathogenesis of a wide range of diseases, including endotoxic shock, cerebral malaria, tumors, autoimmuine disease, AIDS and graft-host rejection (Beutler, B. and Von Huffel, C., Science 264:667-668 (1994)). Mutations in the p55 Receptor cause increased susceptibility to microbial infection.\nMoreover, an about 80 amino acid domain near the C-terminus of TNFR1 (P55) and Fas was reported as the xe2x80x9cdeath domain,xe2x80x9d which is responsible for transducing signals for programmed cell death (Tartaglia et al., Cell 74:845 (1993)). Other regions of the TNF receptor intracellular domain are responsible for the activation of transcription through NF-kB (Cheng and Baltimore Genes and Development 10: 963-973 (1996)). More recent evidence has suggested that receptors may induce signals in cells expressing membrane bound TNF family ligand in a process known as xe2x80x9creverse signalingxe2x80x9d (Wiley et al., J. Immunol. 157:3635-3639.\nThe effects of TNF family ligands and TNF family receptors are varied and influence numerous functions, both normal and abnormal, in the biological processes of the mammalian system. There is a clear need, therefore, for identification and characterization of such receptors and ligands that influence biological activity, both normally and in disease states. In particular, there is a need to isolate and characterize novel members of the TNF receptor family.\nThis indicates that these Tumor necrosis factor receptors (TNF-R) have an established, proven history as therapeutic targets. Clearly there is a need for identification and characterization of further members of Tumor necrosis factor receptor (TNF-R) family which can play a role in preventing, ameliorating or correcting dysfunctions or diseases, including, but not limited to, chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (eg inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, AIDS, Bone diseases, cancer (eg lymphoproliferative disorders), atheroschlerosis, and Alzheimers disease.\nThe present invention relates to tumor necrosis factor receptor (TNF-R) related polypeptides and their ligands, hereinafter referred to as TR1, TR2, TL2 and TL4. The invention also relates to methods to identify agonists and antagonists of TR1, TR2, TL2 and TL4. The agonists and antagonists thus identified can be used to treat chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (eg inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, AIDS, bone diseases, cancer (eg lymphoproliferative disorders), atheroschlerosis, and Alzheimers disease, among others, caused by imbalance of TR1, TR2, TL2 or TL4.\nxe2x80x9cTR1 or TR1 polypeptide or TR1 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:1 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 1; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:1 over its entire length. Furthermore, TR1 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 1. xe2x80x9cTR1 or TR1 polypeptide or TR1 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of the above polypeptides, and some of these derivatives are further illustrated below.\nxe2x80x9cTR2 or TR2 polypeptide or TR2 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:2 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 2; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:2 over its entire length. Furthermore, TR2 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 2. xe2x80x9cTR2 or TR2 polypeptide or TR2 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of the above polypeptides, and some of these derivatives are further illustrated below.\nxe2x80x9cTL2 or TL2 polypeptide or TL2 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:3 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 3; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:3 over its entire length. Furthermore, TL2 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 3. xe2x80x9cTL2 or TL2 polypeptide or TL2 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of the above polypeptides, and some of these derivatives are further illustrated below.\nxe2x80x9cTL4 or TL4 polypeptide or TL4 proteinxe2x80x9d refers, among others, generally to a polypeptide having the amino acid sequence set forth in SEQ ID NO:4 as well as polypeptides comprising the amino acid sequence of SEQ ID NO: 4; and polypeptides comprising the amino acid sequence which have at least 70% identity to that of SEQ ID NO:4 over its entire length. Furthermore, TL4 also refers to a polypeptide which comprises a sequence which has 70% identity to a fragment of SEQ ID NO: 4. xe2x80x9cTL4 or TL4 polypeptide or TL4 proteinxe2x80x9d also includes derivatives, such as fusion proteins, of above polypeptides, and some of these derivatives are further illustrated below.\ncDNA encoding polypeptide of SEQ ID NO:4 is contained in SEQ ID NO:5.\nxe2x80x9cPolypeptidexe2x80x9d refers to any peptide or protein comprising two or more amino acids joined to each other by peptide bonds or modified peptide bonds, i.e., peptide isosteres. xe2x80x9cPolypeptidexe2x80x9d refers to both short chains, commonly referred to as peptides, oligopeptides or oligomers, and to longer chains, generally referred to as proteins. Polypeptides may contain amino acids other than the 20 gene-encoded amino acids. xe2x80x9cPolypeptidesxe2x80x9d include amino acid sequences modified either by natural processes, such as posttranslational processing, or by chemical modification techniques which are well known in the art. Such modifications are well described in basic texts and in more detailed monographs, as well as in a voluminous research literature. Modifications can occur anywhere in a polypeptide, including the peptide backbone, the amino acid side-chains and the amino or carboxyl termini. It will be appreciated that the same type of modification may be present in the same or varying degrees at several sites in a given polypeptide. Also, a given polypeptide may contain many types of modifications. Polypeptides may be branched as a result of ubiquitination, and they may be cyclic, with or without branching. Cyclic, branched and branched cyclic polypeptides may result from posttranslation natural processes or may be made by synthetic methods. Modifications include acetylation, acylation, ADP-ribosylation, amidation, covalent attachment of flavin, covalent attachment of a heme moiety, covalent attachment of a nucleotide or nucleotide derivative, covalent attachment of a lipid or lipid derivative, covalent attachment of phosphotidylinositol, cross-linking, cyclization, disulfide bond formation, demethylation, formation of covalent cross-links, formation of cystine, formation of pyroglutamate, formylation, gamma-carboxylation, glycosylation, GPI anchor formation, hydroxylation, iodination, methylation, myristoylation, oxidation, proteolytic processing, phosphorylation, prenylation, racemization, selenoylation, sulfation, transfer-RNA mediated addition of amino acids to proteins such as arginylation, and ubiquitination. See, for instance, PROTEINSxe2x80x94STRUCTURE AND MOLECULAR PROPERTIES, 2nd Ed., T. E. Creighton, W. H. Freeman and Company, New York, 1993 and Wold, F., Posttranslational Protein Modifications: Perspectives and Prospects, pgs. 1-12 in POSTTRANSLATIONAL COVALENT MODIFICATION OF PROTEINS, B. C. Johnson, Ed., Academic Press, New York, 1983; Seifter et al., xe2x80x9cAnalysis for protein modifications and nonprotein cofactorsxe2x80x9d, Meth Enzymol (1990) 182:626-646 and Rattan et al., xe2x80x9cProtein Synthesis: Posttranslational Modifications and Agingxe2x80x9d, Ann NY Acad Sci (1992) 663:48-62.\nxe2x80x9cIdentityxe2x80x9d is a measure of the identity of nucleotide sequences or amino acid sequences. In general, the sequences are aligned so that the highest order match is obtained. xe2x80x9cIdentityxe2x80x9d per se has an art-recognized meaning and can be calculated using published techniques. See, e.g.: (COMPUTATIONAL MOLECULAR BIOLOGY, Lesk, A. M., ed., Oxford University Press, New York, 1988; BIOCOMPUTING: INFORMATICS AND GENOME PROJECTS, Smith, D. W., ed., Academic Press, New York, 1993; COMPUTER ANALYSIS OF SEQUENCE DATA, PART I, Griffin, A. M., and Griffin, H. G., eds., Humana Press, New Jersey, 1994; SEQUENCE ANALYSIS IN MOLECULAR BIOLOGY, von Heinje, G., Academic Press, 1987; and SEQUENCE ANALYSIS PRIMER, Gribskov, M. and Devereux, J., eds., M Stockton Press, New York, 1991). While there exist a number of methods to measure identity between two polynucleotide or polypeptide sequences, the term xe2x80x9cidentityxe2x80x9d is well known to skilled artisans (Carillo, H., and Lipton, D., SIAM J Applied Math (1988) 48:1073). Methods commonly employed to determine identity or similarity between two sequences include, but are not limited to, those disclosed in Guide to Huge Computers, Martin J. Bishop, ed., Academic Press, San Diego, 1994, and Carillo, H., and Lipton, D., SIAM J Applied Math (1988) 48:1073. Methods to determine identity and similarity are codified in computer programs. Preferred computer program methods to determine identity and similarity between two sequences include, but are not limited to, GCG program package (Devereux, J., et al., Nucleic Acids Research (1984) 12(1):387), BLASTP, BLASTN, FASTA (Atschul, S. F. et al., J Molec Biol (1990) 215:403).\nAs an illustration, by a polypeptide having an amino acid sequence having at least, for example, 70% xe2x80x9cidentityxe2x80x9d to a reference amino acid sequence of SEQ ID NO:2 is intended that the amino acid sequence of the polypeptide is identical to the reference sequence except that the polypeptide sequence may include up to thirty amino acid alterations per each 100 amino acids of the reference amino acid of SEQ ID NO: 2. In other words, to obtain a polypeptide having an amino acid sequence at least 70% identical to a reference amino acid sequence, up to 30% of the amino acid residues in the reference sequence may be deleted or substituted with another amino acid, or a number of amino acids up to 30% of the total amino acid residues in the reference sequence may be inserted into the reference sequence. These alterations of the reference sequence may occur at the amino or carboxy terminal positions of the reference amino acid sequence or anywhere between those terminal positions, interspersed either individually among residues in the reference sequence or in one or more contiguous groups within the reference sequence.\nWe have now discovered that TL2 of SEQ ID NO: 3 (otherwise known as TRAIL or Apo-2L, (Wiley S R, et al., Immunity (6):673-682 (1995); Pitt et al., J. Biol. Chem. 271: 12687-12690 (1996)) is a ligand of TR1 of SEQ ID NO:1 (otherwise known as osteoprotegerin (OPG), W. S. Simonet, et al., Cell, Vol 89, pp 309-319, 1997). Further, we also discovered that TL4 of SEQ ID NO: 4 is a ligand of TR2 of SEQ ID NO: 2 (described by R. I. Montgomery et al., Cell, Vol 87, pp427-436, 1996; Kwon et al., J. Biol. Chem. 272: 14272-14276 (1997); Hsu et al., J. Biol. Chem. 272:13471-13474 (1997))). Thus, the TR1 and TR2 polypeptides of the present invention, and their respective ligands, TL2 and TL4, can be employed in a screening process for compounds which bind to the receptors, or to their ligands, and which activate (agonists) or inhibit activation of (antagonists) TR1 and TR2 receptor polypeptides of the present invention, or their respective ligands TL2 and TL4. Thus, polypeptides of the invention may be used to assess the binding of small molecule substrates, receptors and ligands in, for example, cells, cell-free preparations, chemical libraries, and natural product mixtures. These substrates, receptors and ligands may be natural substrates and ligands, or may be structural or functional mimetics. See Coligan et al., Current Protocols in Immunology 1(2):Chapter 5 (1991).\nTR1, TR2, TL2 and TL4 polypeptides are responsible for many biological functions, including many pathologies. Accordingly, it is desirous to find compounds and drugs which stimulate TR1, TR2, TL2 or TL4 on the one hand, and which can inhibit the function of TR1, TR2, TL2 or TL4 or remove TR1, TR2, TL2 or TL4 expressing cells on the other hand (also defined as antagonists). Antagonists for TR1, TR2, TL2, and TL4 (including agents which remove TR1, TR2, TL2 or TL4 expressing cells) may be employed for a variety of therapeutic and prophylactic purposes for such conditions as chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (e.g. inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, AIDS, Bone diseases, cancer (e.g. lymphoproliferative disorders), atheroschlerosis, and Alzheimers disease. Agonists can be employed for therapeutic and prophylactic purposes for such conditions responsive to activation of T cells and other components of the immune system, such as for treatment of cancer and AIDS. However, agonists can also be employed for inappropriate stimulation of T cells and other components of the immune system which leads to down modulation of immune activity with therapeutic or prophylactic application for conditions such, as chronic and acute inflammation, arthritis, septicemia, autoimmune diseases (e.g. inflammatory bowel disease, psoriasis), transplant rejection, graft vs. host disease, infection, stroke, ischemia, acute respiratory disease syndrome, restenosis, brain injury, bone diseases, atheroschlerosis, and Alzheimers disease.\nAntagonists may be identified using assays to detect compounds which inhibit binding of TL2 to TR1 (or TL4 to TR2) in either cell-free or cell based assays. Suitable cell-free assays may be readily determined by one of skill in the art. For example, an ELISA format may be used in which purified TR1 (or TR2), or a purified derivative of TR1 (or TR2), such as a fusion protein, containing the extracellular domain of TR1 (or TR2), is immobilized on a suitable surface, either directly or indirectly (e.g., via an antibody to TR1 (or TR2) or to the fused epitope or protein domain) and candidate compounds are identified by their ability to block binding of purified soluble, extracellular domain of TL2 to TR1 (or soluble, extracellular domain of TL4 to TR2). The binding of TL2 to TR1 (or TL4 to TR2) could be detected by using a label directly or indirectly associated with TL2 (or TL4). Suitable detection systems include the streptavidin horseradish peroxidase conjugate, or direct conjugation by a tag, e.g., fluorescein. Conversely, purified, soluble TL2 (or TL4) may be immobilized on a suitable surface, and candidate compounds identified by their ability to block binding of purified TR1 to TL2 (or TR2 to TL4). The binding of TR1 to TL2 (or TR2 to TL4) could be detected by using a label directly or indirectly associated with TR1 (or TR2). Many other assay formats are possible that use the TR1 (or TR2) protein and its ligands.\nSuitable cell based assays may be readily determined by one of skill in the art. In general, such screening procedures involve producing appropriate cells which express the receptor polypeptides (or ligands thereof) of the present invention on the surface thereof. Such cells include cells from mammals, yeast, Drosophila or E. coli. Cells expressing the receptor, such as TR1 or TR2, (or cell membrane containing the expressed receptor) are then contacted with a ligand, such as TL2 or TL4, or test compound to observe binding, or stimulation or inhibition of a functional response. The assays may simply test binding of a candidate compound wherein adherence to the cells bearing the receptor is detected by means of a label directly or indirectly associated with the candidate compound or in an assay involving competition with a labeled competitor, such as the ligands TL2 or TL4. Alternatively, cells expressing the ligand, such as TL2 or TL4, (or cell membrane containing the expressed ligand) are then contacted with a receptor, such as TR1 or TR2, or test compound to observe binding, or stimulation or inhibition of a functional response. Similarly, the assays may simply test binding of a candidate compound wherein adherence to the cells bearing the ligand is detected by means of a label directly or indirectly associated with the candidate compound or in an assay involving competition with a labeled competitor, such as the receptors TR1 and TR2. Further, these assays may test whether the candidate compound results in a signal generated by activation of the receptor (e.g TR1 or TR2) or its respective ligand (e.g. TL2 or TL4) using detection systems appropriate to the cells bearing the receptor or its ligand and fusion proteins thereof at their surfaces. Typical fusion partners include fusing the extracellular domain of the receptor or ligand with the intracellular tyrosine kinase domain of a second receptor. Inhibitors of activation are generally assayed in the presence of an agonist, such as the ligand TL2 or TL4 for cells expressing TR1 and TR2 receptors or receptor fusions respectively, or the receptor TR1 and TR2 with cells expressing TL2 and TL4 ligands and ligand fusions, and the effect on activation by the agonist by the presence of the candidate compound is observed. Standard methods for conducting such screening assays are well understood in the art.\nExamples of potential TR1 or TR2 antagonists include antibodies or, in some cases, proteins which are closely related to the ligand of the TR1 or TR2, e.g., a fragment of the respective ligand TL2 or TL4, or small molecules which bind to the receptor, or its ligand, but do not elicit a response, so that the activity of the receptor is prevented. Examples of potential TR1 or TR2 agonists include antibodies that bind to TR1 or TR2, its respective ligand, such as TL2 or TL4, or derivatives thereof, and small molecules that bind to TR1 or TR2. These agonists will elicit a response mimicking all or part of the response induced by contacting the native ligand.\nSince receptors may also invoke signals in cells expressing the membrane TL2 and TL4, these screens may also yield agonists which mimic the agonist activity of TR1 with membrane TL2 and TR2 with membrane TL4. Examples of potential TL2 or TL4 agonists include antibodies that bind to TL2 or TL4, its respective receptor, such as TR1 or TR2, or derivatives thereof, and small molecules that bind to TL2 or TL4. These agonists will elicit a response mimicking all or part of the response induced by contacting the native ligand. Alternatively, TR1 or TR2 may be expressed as a soluble protein, including versions which fuse all or part of TR1 or TR2 with a convenient partner peptide for which detection reagents are available, eg TR1-IgG or TR2-IgG fusions, and used in a solid state or solution phase binding assay. For example, the soluble TR1 or TR2 can be used to detect agonist or antagonist binding directly through changes that can be detected experimentally, eg surface plasmon resonance, nuclear magnetic resonance spectrometry, sedimentation, calorimetry. The soluble TR1 or TR2 can be used to detect agonist or antagonist binding indirectly by looking for competition of the candidate agonist or antagonist with a ligand, such as TL2 or TL4, whose binding can be detected. Ligand detection methods include antibody recognition, modification of the ligand via radioactive labeling, chemical modification (eg biotinylation), fusion to an epitope tag. Methods include ELISA based assays, immunoprecipitation and scintillation proximity.\nAssays similar to those described above using soluble or membrane bound TR1 or TR2 may also be used to identify and purify additional natural ligand(s) of TR1 or TR2. These ligands may be agonists or antagonists of the receptor.\nThus the invention relates to:\nI. A method for identifying agonists or antagonists to TR1 or TR2 comprising:\n(a) contacting a candidate compound with TR1 or TR2 in the presence of labeled or unlabeled ligand TL2 or TL4 respectively; and\n(b) assessing the ability of said candidate compound to compete with TL2 or TL4 binding to TR1 or TR2 respectively;\nII. The method of I in which TR1 or TR2 is on the surface of a host cell, on a cell membrane or on a solid support;\nIII. The method of II for identifying agonists which further includes determining whether the candidate compound affects a signal generated by TR1 or TR2 polypeptide at the surface of the cell, wherein a candidate compound which increases production of said signal is identified as an agonist;\nIV. An agonist identified by the method of I, II or III;\nV. The method of II for identifying antagonists which further includes determining whether the candidate compound affects a signal generated by TR1 or TR2 polypeptide at the surface of the cell, wherein a candidate compound which diminishes production of said signal is identified as an antagonist;\nVI. An antagonist identified by the method of I, II or V;\nVII. A method for identifying agonists or antagonists to TL2 or TL4 comprising:\n(a) contacting a candidate compound with TL2 or TL4 in the presence of labeled or unlabeled TR1 or TR2 respectively; and\n(b) assessing the ability of said candidate compound to compete with TR1 or TR2 binding to TL2 or TL4 respectively;\nVIII. The method of VII in which TL2 or TL4 is on the surface of a host cell, on a cell membrane or on a solid support;\nIX. The method of VIII for identifying agonists to TL2 or TL4 which includes determining whether the candidate compound affects a signal generated by TL2 or TL4 polypeptide at the surface of the cell, wherein a candidate compound which increases production of said signal is identified as an agonist;\nX. An agonist identified by the method of IX;\nXI. The method of VIII for identifying antagonists which further includes determining whether the candidate compound affects a signal generated by TL2 or TL4 polypeptide at the surface of the cell, wherein a candidate compound which diminishes production of said signal is identified as an antagonist; and\nXII. An antagonist identified by the method XI.\nThus in another aspect, the present invention relates to a screening kit for identifying agonists, antagonists, ligands, receptors, substrates, enzymes, etc. for TR1, TR2, TL2 or TL4 polypeptides; which comprises:\n(a) a TR1, TR2, TL2 or TL4 polypeptide, preferably that of SEQ ID NO:1, 2, 3 or 4;\n(b) a recombinant cell expressing a TR1, TR2, TL2 or TL4 polypeptide, preferably that of SEQ ID NO: 1, 2, 3or 4;\n(c) a cell membrane expressing a TR1, TR2, TL2 or TL4 polypeptide; preferably that of SEQ ID NO: 1,2,3 or4; or\n(d) antibody to a TR1, TR2, TL2 or TL4 polypeptide, preferably that of SEQ ID NO: 1, 2, 3 or 4. It will be appreciated that in any such kit, (a), (b), (c) or (d) may comprise a substantial component.\nProphylactic and Therapeutic Methods\nThis invention provides methods of treating an abnormal conditions related to both an excess of and insufficient amounts of TR1, TR2, TL2 or TL4 polypeptide activity.\nIf the activity of TR1, TR2, TL2 or TL4 polypeptide is in excess, several approaches are available. One approach comprises administering to a subject an inhibitor compound (antagonist) as hereinabove described along with a pharmaceutically acceptable carrier in an amount effective to inhibit activation by blocking binding of ligands to the TR1 or TR2 polypeptide, or by inhibiting a second signal, and thereby alleviating the abnormal condition.\nIn another approach, soluble forms of TR1 or TR2 polypeptides still capable of binding the ligand in competition with endogenous TR1 or TR2 polypeptide may be administered. Typical embodiments of such competitors comprise fragments of the TR1 or TR2 polypeptide.\nFor treating abnormal conditions related to an under-expression of TR1, TR2, TL2 or TL4 and its activity, several approaches are also available. One approach comprises administering to a subject a therapeutically effective amount of a compound which activates TR1, TR2, TL2 or TL4 polypeptide, i.e., an agonist as described above, in combination with a pharmaceutically acceptable carrier, to thereby alleviate the abnormal condition.\nFormulation and Administration\nAgonists and antagonist of TR1, TR2, TL2 or TL4 may be formulated in combination with a suitable pharmaceutical carrier. Such formulations comprise a therapeutically effective amount of the polypeptide or compound, and a pharmaceutically acceptable carrier or excipient. Such carriers include but are not limited to, saline, buffered saline, dextrose, water, glycerol, ethanol, and combinations thereof. Formulation should suit the mode of administration, and is well within the skill of the art. The invention further relates to pharmaceutical packs and kits comprising one or more containers filled with one or more of the ingredients of the aforementioned compositions of the invention.\nPolypeptides and other compounds of the present invention may be employed alone or in conjunction with other compounds, such as therapeutic compounds.\nPreferred forms of systemic administration of the pharmaceutical compositions include injection, typically by intravenous injection. Other injection routes, such as subcutaneous, intramuscular, or intraperitoneal, can be used. Alternative means for systemic administration include transmucosal and transdermal administration using penetrants such as bile salts or fusidic acids or other detergents. In addition, if properly formulated in enteric or encapsulated formulations, oral administration may also be possible. Administration of these compounds may also be topical and/or localized, in the form of salves, pastes, gels and the like.\nThe dosage range required depends on the choice of peptide or compound, the route of administration, the nature of the formulation, the nature of the subject\"\"s condition, and the judgment of the attending practitioner. Suitable dosages, however, are in the range of 0.1-100 xcexcg/kg of subject. Wide variations in the needed dosage, however, are to be expected in view of the variety of compounds available and the differing efficiencies of various routes of administration. For example, oral administration would be expected to require higher dosages than administration by intravenous injection. Variations in these dosage levels can be adjusted using standard empirical routines for optimization, as is well understood in the art."}
{"text": "The present invention relates to the production of magnesium metal by the reduction of magnesium oxide at elevated temperatures in the presence of an aluminum metal reducing agent and a molten oxidic slag, in an electric furnace, and the condensation of vaporized magnesium in a condenser.\nAn advantageous method of producing magnesium lies in the chemical reduction of magnesium oxide with a reducing agent in the presence of a molten oxide slag, in an electric furnace. A variety of metallothermic processes, as they are commonly called, for the production of magnesium employing various reducing agents, various types of reactants, and under various conditions of temperatures and pressures have been proposed.\nIn general, the various metallothermic processes are concerned with the production of metallic magnesium by the reduction of magnesium oxide with a metal reductant at elevated temperatures. Magnesium oxide, usually in the form of calcined dolomite (dolime) or calcined magnesite or mixtures thereof, is caused to react with a metallic reducing agent, such as silicon, aluminum, calcium or mixtures or alloys thereof, in the presence of a molten slag bath in a furnace at temperatures in excess of 1300.degree. C, to release magnesium vapor which may be condensed and collected. Some of the processes are carried out in the presence of an inert gas.\nAn early process of this type, called the Pidgeon process provides for the production of magnesium by the reduction of magnesium oxide with ferrosilicon.\nA more recent process is exemplified by U.S. Pat. 2,971,833 and is known as the Magnetherm process. The process is operated under a very high vacuum (pressure ranging between 5 and 20 millimeters of mercury) and at a temperature of about 1500.degree. C utilizing an electric furnace. Silicon is employed as the reductant, preferably in the form of ferrosilicon containing 70-80 percent Si, silicon of a purity in excess of 97%, or an aluminum-ferrosilicon. Care is taken to avoid that the silicon content of the residual ferrosilicon drops below 33.5%. A magnesia containing substance is dissolved in a liquid slag consisting essentially of lime, silica and alumina wherein the ratio of calcium oxide to silicon dioxide and the ratio of aluminum oxide to silicon dioxide are controlled.\nProcesses employing aluminum metal as the reductant are exemplified by U.S. Pat. Nos. 3,782,922; 2,527,722 and 2,527,724. The first patent discloses the production of magnesium by the reduction of magnesium oxide from a mixture of magnesium oxide and calcium oxide with a reducing agent comprising at least 85% aluminum in the presence of a molten calcium-aluminate slag bath at a temperature of about 1300.degree.-1700.degree. C and a pressure of about atmospheric. Slag is removed when the magnesium oxide content thereof is less than five percent, with the slag comprising 35-65 percent aluminum oxide, 35-55 percent calcium oxide and 0-10 percent silicon dioxide. The other two patents utilize magnesium silicate in pure or ore form and aluminum as reactants. Dolomite is also used in the former and the magnesium silicate is of a particular granular size in the latter.\nU.S. Pat. No. 3,658,509 also discloses the use of aluminum as the reductant, as well as silicon and aluminum-silicon alloys with the latter being preferred. In such process, an inert gas is used to obviate at least in part the need of a high vacuum, and the slag contains 20-50 percent silicon dioxide.\nCopending application Ser. No. 610,430, now U.S. Pat. No. 4,033,759 discloses the use of aluminum metal as a reductant in such a manner as to produce a slag wherein, by weight percent, calcium oxide ranges from 30-65, aluminum oxide ranges from 28-64, magnesium oxide ranges from 6-13 and silicon dioxide is equal to or less than 5.\nThe primary object of this invention is to provide an improved process for production of magnesium metal wherein aluminum metal is reacted with a calcium magnesium aluminate slag, which eliminates the disadvantages of the prior art processes, while retaining the benefits thereof.\nAnother object of this invention is to provide a process for production of magnesium metal at atmospheric pressure by reaction of aluminum metal with a calcium magnesium aluminate slag, or with magnesium oxide in the presence of such slag, wherein the silicon dioxide content of the slag is maintained at a low level.\nAnother object of the invention is to provide a process for producing magnesium metal utilizing aluminum metal as the reductant, wherein magnesium vapor is produced at a pressure of one atmosphere thereby eliminating the need for vacuum equipment and permitting a continuous operation at such normal pressure.\nAn important object of the present invention is to provide a process of the foregoing type which is economical and which can utilize aluminum metal scrap as the metal reductant.\nOther objects and advantages of the present invention will become more readily apparent from a consideration of the description hereinafter."}
{"text": "A need exists in the building industry for prefabricated building structures which may be erected quickly with minimal site preparation by people with relatively little job skills. Such structures should be able to withstand the natural elements, provide a large degree of usable interior space and be aesthetically pleasing to the eye. Prefabrication of building components can assure quality and standardization of parts by providing a factory controlled environment for production. Production losses due to adverse weather conditions can be minimized. Greatly increased efficiency results and aids in the reduction of construction costs.\nThe prefabricated building structure of the present invention uses a polygonal or \"round\" design. The \"round\" configuration of the present invention provides numerous advantages. One such advantage is that of providing a specified amount of interior floor space within less outside wall space (perimeter feet) than a comparable square or rectangular building, because the circumference of a circle is less than the perimeter of a rectangle for a given area. For example, a round building with 1000 square feet has approximately 112 perimeter feet. A rectangular building with 1000 square feet could be 20'.times.50' having a perimeter of 140 feet. Since heat loss is directly proportional to the number of perimeter feet (outside wall space), heat loss in the present invention is greatly reduced. Other advantages include minimized interior obstructions and increased outside wall strength inherent to a curved or arch segment.\nThe structure of the present invention is extremely flexible as to usage because of the minimized interior obstructions. The structure has only one interior central support column. The finished structure may have residential, commercial, recreational, or educational use. The structure may be totally enclosed or partially opened. Walls may be constructed from conventional building materials or may utilize large expanses of glass or solar type glass panels.\nThe roof and floor components of the invention utilize smaller lumber sizes than conventional construction which has the effect of utilizing random lumber lengths, thus conserving lumber and reducing material costs. Further, the roof and floor components possess a triangle-shape symmetry which enables them to be easily stacked or nested for storage and shipment.\nThe vertically adjustable perimeter columns of the present invention allow for minimal site preparation on uneven terrain. Conventional poured footings, foundations and masonry walls are all eliminated. Reduced construction time and reduced costs are a result.\nThe building structure of the present invention requires no skilled masons, concrete finishers, or carpenters for construction. The structure may be constructed, using six men with minimal labor skills, in approximately eight hours."}
{"text": "The present invention relates to an adhesive, the use thereof, a method for producing the same, an engineered wood board, and a method for producing the engineered wood board.\nAdhesive are a decisive constituent in the industrial production of a multitude of products, for example in the production of wood-based materials. Adhesives can be classified according to different criteria based on the chemical basis of the adhesives or their setting mechanism. Proceeding from the production method of the adhesives, a total of three superordinate classes of adhesives are defined: Polymerization adhesives, polyaddition adhesives, and polycondensation adhesives. Further divisions of the adhesives with respect to their physical and/or chemical properties, such as e.g. hot-melt adhesives or solvent adhesives, also are possible.\nThe polymerization adhesives are produced by reaction of monomers including a carbon-carbon double bond after activation. The activation of the starting materials can be effected by suitable catalysts or radicals and/or in the presence of radiation, for example UV radiation or electron beam. Typical polymerization adhesives are to be assigned e.g. to the group of acrylate adhesives.\nA polycondensation adhesive, on the other hand, can be obtained by reaction of two monomer molecules by splitting off a simple molecule, such as water, acid or alcohol. The polymeric reaction product thus is present together with a by-product obtained during the reaction, so that corresponding measures are required during the processing of these adhesives. The most important polycondensates for use as adhesives include polyamides, polyesters and silicones or formaldehyde condensates, wherein reference should be made here in particular to the phenol-formaldehyde adhesives (PF), cresol/resorcinol-formaldehyde adhesives, urea-formaldehyde-resin adhesives (UF) or melamine-formaldehyde adhesives (MF).\nThe production of the polyaddition adhesive is based on the addition of various reactive monomer molecules by simultaneous migration of a hydrogen atom from the one component to the other component. Typical representatives include epoxy resin adhesives, polyurethanes or polycyanurates.\nFor the production of wood-based materials such as engineered wood boards, wood chipping products are coated with the adhesive and compressed to molded articles by applying pressure and temperature.\nThe type of adhesive used is substantially influenced by the size and quality of the wood fibers and/or wood chips used.\nFor example in the production of wood particle boards and wood fiber boards, such as e.g. MDF and HDF boards, which are produced from wood fibers in a dry process, there are frequently used polycondensation adhesives, in particular in the form of urea-formaldehyde resins. The particular advantage in the use of formaldehyde resins as adhesives consists in their high availability, the low costs as well as an easy manufacturability and handleability. Since the formaldehyde resins usually are produced by reaction with an excess of formaldehyde, these excesses of formaldehyde also are detectable in the intermediate and/or end products. Since formaldehyde, however, is classified as carcinogenic, the use of formaldehyde resins thus turns out to be disadvantageous in particular for the production of engineered wood boards for use indoors. In addition, the resins condensed-out have a low water stability.\nIn the production of boards from oriented wood chips, so-called OSB boards, on the other hand, polyaddition adhesives containing urethanes, e.g. on the basis of diphenyl methane diisocyanate adhesives, are used more and more. The complete chemical reaction without disturbing excesses and the high adhesive force are regarded as particular advantages of polyurethane adhesives. What is regarded as particularly disadvantageous, on the other hand, is the limited availability at a high price and the affinity to metal, so that metallic tools and plant sections must particularly be protected against direct contact. It is also disadvantageous that e.g. PMDI reacts already with the water from the room air humidity.\nIt is known that with a joint use of polycondensation adhesives and polyaddition adhesives incompatibilities will occur between the two adhesive systems, which will lead to comparatively poor technological properties. This problem is particularly pronounced at the interfaces of layers with different glue systems."}
{"text": "The present disclosure relates generally to devices for measuring fluid pressure, and more particularly to devices for invasively measuring fluid pressure.\nThe monitoring of blood pressure in patients undergoing hemodialysis or infusion is usually performed with a non-disposable device. Such devices may include electronic sensors that have a saline buffer between the fluid and the sensor. One potential problem with non-disposable fluid pressure measuring devices is that the spread of disease may be substantially difficult to control. Various fluid-pressure devices use fluid (liquid or gel), generally contained between two diaphragms, as the media to transmit the pressure to the sensor that monitors pressure. One of the diaphragms conforms responsive to the pressure received, while the other contains the fluid.\nMany diaphragms used in hemodialysis or infusion blood-pressure devices are in an elastic condition, which generally requires calibration between uses. These devices generally use both: well-controlled materials for the diaphragm and/or mounting of the diaphragm; and a calibration procedure that accounts for mounting and material variability. Such materials and calibration generally undesirably add to the expense of the device and efficiency of use thereof.\nAs such, it would be desirable to provide a device that is disposable, thereby substantially aiding in the prevention of bio-contamination from one patient to the next. Further, it would be desirable to provide such a device that generally does not require calibration between uses."}
{"text": "Ladders are conventionally utilized to provide a user thereof with improved access to elevated locations that might otherwise be inaccessible. Ladders come in many shapes and sizes, such as straight ladders, straight extension ladders, stepladders, and combination step and extension ladders. So-called combination ladders may incorporate, in a single ladder, many of the benefits of various ladder designs.\nLadders known as step ladders, sometimes referred to as A-frame ladders, are self-supporting ladders, meaning that they do not need to be leaned against a wall, pole or other structure for stability. Rather, step ladders may be positioned on a floor (or other similar surface) such that at least three feet of the ladder, and conventionally four feet, provide a stable support structure for a user to climb upon, even in an open space (e.g., outside or in the middle of a room) without a wall, roof, pole or other type of structure being necessary for the stability of the ladder. Conventional step ladders may include a first rail assembly coupled with a top cap and a second rail assembly coupled with the top cap. One of the rail assemblies conventionally includes a plurality of rungs that are evenly spaced between the supporting surface (e.g., the floor or ground) and the top cap.\nRegardless of the type of ladder being employed, using a ladder can present various risks to the user. For example, one potential hazard exists where a user is distracted or is not paying attention and loses track of which rung they are currently standing on—particularly when they are descending the ladder. In such a case, a user may think, for example, that their next “step” downward will place them on the ground at the bottom of the ladder when, in reality, there is still one more step for them to descend prior to reaching the ground. This misperception can result in the user stumbling on, or even missing completely, the lowest rung of the ladder. There is a continuing desire in the industry to provide ladders that reduce the risk of accident and provide improved safety and stability to a user thereof."}
{"text": "Emerging mobile devices will provide access to video, audio and other content, and additional content-related information. As only one example, existing mobile information browsing devices that feature content browsing require a continuous two-way connection with the content host server by means of a two-way radio connection, such as by WiFi or the cellular public service telephone network. Such a system constrains real-time content browsing to the physical location of the wireless access point (WAP), or to the use of a cellular network that incurs usage tariffs. The instant invention makes content browsing possible without the need for constant proximity to a WAP, and without incurring cellular network tariffs, or other access costs. More generally, the invention provides a system and method that allows a content provider or service manager to provide different levels of service, as well as to provide for the most efficient utilization of the available communication channels."}
{"text": "The present invention generally relates to illuminated panel portions for vehicles and vehicular accessories. More specifically, the present invention relates to an illuminated panel portion including a generally translucent panel with an outer surface that is substantially flush with the outer surface of the adjacent panel of vehicles and vehicular accessories for illuminating a selected design defined by the translucent panel.\nIlluminated devices for vehicles and vehicular accessories have been traditionally used to enhance the visibility of the vehicle or passenger. Also, illuminated devices have served to illuminate selected designs such as indicia, emblems, pictures, outlines, or other such artistic and other designs.\nA number of illuminated devices have been designed to be mounted on various vehicle or vehicular accessory surfaces. For example, U.S. Pat. No. 6,190,026 describes an illuminated automotive emblem for a motor vehicle having a translucent display template that is illuminated from behind by a light source. U.S. Pat. Nos. 5,615,940 and 5,479,324 describes an illumination device mounted on a gas tank or shell.\nThese and other prior illumination devices are typically mounted onto or protrude from the surfaces of the vehicle or vehicular accessory. Such a mounting arrangement can actually increase the possibility of damage to the underlying machinery or circuitry protected by the support surface. For example, water may seep and pool between the mounting surface of the illumination device and the surface of the vehicle or vehicle accessory. If the illumination device covers an aperture in the vehicle or vehicle accessory surface, water can sometimes seep under the mounting surface of the illumination device and damage underlying machinery or circuitry.\nIllumination devices are also typically mounted or protrude above the surface of the vehicle or vehicular accessory. This mounting arrangement alters the preexisting contour of the surface. By altering the contour of the vehicle or vehicle accessory surface, the aerodynamics, appearance or aesthetics, and performance of the vehicle and/or accessory may also be compromised. These potential drawbacks may be a disincentive for consumers to mount or install such illumination devices on their vehicles and/or vehicular accessories.\nIn a variation of the above mounting techniques U.S. Pat. No. 6,439,589 shows a mounting arrangement in which the desired design of the indicia is cut from the surface of the vehicle or vehicular accessory. This arrangement causes the edges of the cut-out portion to be raised above the outer surface of the illumination device. Moreover, because the indicia is cut out, the complexity of the indicia design is necessarily limited in order to maintain the structural integrity of the support surface of the vehicle or the vehicular accessory.\nTherefore, it is an aspect and object of this invention to provide an illuminated panel portion including a generally translucent panel with an outer surface which avoids one or more of the drawbacks of prior illumination devices."}
{"text": "This invention relates to sense amplifier circuits used in memory circuits, and more particularly, this invention relates to a sense amplifier circuits having a global data-bus line and reduced power consumption and noise.\nSense amplifiers are often used in static random access memories (SCRAM) and other types of memories during the Read operation of data from a memory cell. The memory cells are etched onto silicon wafers in bit lines formed of columns and word lines formed as rows. The intersection of a bit line and word line form an address of the memory cell. During a write operation, a row line contains a charge that the capacitor/transistor should take to place a value at the address corresponding to a xe2x80x9c1xe2x80x9d or xe2x80x9c0xe2x80x9d. When reading, a sense amplifier determines the level of charge of the particular address. For example, if the charge is greater than 50% of a normal charge, the sense amplifier reads a xe2x80x9c1.xe2x80x9d Otherwise, the sense amplifier reads the charge as a xe2x80x9c0xe2x80x9d. A controller and counter tracks a refresh sequence based on the rows that have been accessed in a particular order.\nIn one type of sense amplifier circuit, the differential voltage produced by a memory cell is amplified by a sense amplifier and fed through a single line output driver, for example, a global data-bus line (GDB) such as disclosed in commonly assigned U.S. Pat. No. 5,619,466, the disclosure which is hereby incorporated by reference in its entirety.\nIn that type of sense amplifier circuit, a sense amplifier is coupled to a memory cell via data lines and amplifies the data. The circuit includes a read bus complement (RBC) and read bus true (RBT) line. An equalizer circuit is coupled to the sense amplifier and operable to receive an equalization signal to equalize the sense amplifier. An enable circuit is coupled to the sense amplifier and operable to receive an enable signal to enable the sense amplifier to amplify data.\nDuring a read operation, the read bus complement and read bus true lines are equalized by a sense amplifier equalization (SAEQB) signal that is generated by an internal equalization control circuit during an equalization phase of operation. During this equalization phase, a sense amplifier enable (SAEN) signal is generated by an enable circuit at a xe2x80x9clowxe2x80x9d value such that the cross-coupled transistors forming the sense amplifier are not activated. After the read bus complement and read bus true lines are properly equalized, the sense amplifier equalization (SAEQB) line turns xe2x80x9chighxe2x80x9d corresponding to Vcc. One of the read buses, i.e., the read bus complement or read bus true, will be pulled towards ground or xe2x80x9clowxe2x80x9d by the memory cell. After proper development of the differential in the read bus complement and read bus true line, the sense equalization enable signal goes xe2x80x9chigh,xe2x80x9d turning xe2x80x9cONxe2x80x9d the cross-coupled transistor structure forming the sense amplifier and further pulling down the particular read bus and increasing the voltage differential between the read bus complement and read bus true line, reducing access time. This can create an extra and unnecessary transition during the equalization phase, increasing the power consumption of the overall sense amplifier circuit.\nIt is therefore an object of the present invention to control the data output of a sense amplifier circuit as described above during the equalization phase to reduce unnecessary transitions at the data output.\nThe present invention advantageously provides a sense amplifier circuit for a memory cell that includes a sense amplifier that is operable to be coupled to a memory cell via data lines and operable to amplify data. It includes read bus complement and read bus true lines and a data output for outputting a data output signal. An equalization circuit is coupled to the sense amplifier and operable to receive an equalization signal to equalize the sense amplifier. An enable circuit is coupled to the sense amplifier and operable to receive an enable signal to enable the sense amplifier to amplify data. A control circuit is operable for disconnecting the data output from one of the read bus complement or read bus true lines and minimize unwanted transitions on the data output signal.\nIn one aspect of the present invention, the control circuit comprises a latch circuit that stores data when an enable signal comprises a low signal. It could be formed by at least one passgate circuit that is turned on and off by an enable signal. In yet another aspect of the present invention, the latch circuit is formed as first and second passgate circuits and first and second feedback circuits. The first latch circuit is off when the second latch circuit is on. The control circuit is operable from the enable circuit, in yet another aspect of the present invention. The control circuit can be operable for disconnecting the data output from the read bus complement line whenever there is an equalization. The data output preferably comprises a single line data output such as a global data-bus (GDB) line. The sense amplifier is formed as a pair of cross-coupled transistors of the type known to those skilled in the art.\nA method aspect of the invention controls a sense amplifier circuit and comprises the steps of equalizing a sense amplifier with an equalization signal and disconnecting the data output from one of the read bus complement or read bus true lines during the equalization phase of the sense amplifier. An enabling signal is provided for enabling the sense amplifier. The data output can be disconnected from the read bus complement line. In yet another aspect of the present invention, the data output is disconnected from a control circuit that is operative with a line drive circuit through which the data output signal is passed. The data output can comprise a global data-bus line."}
{"text": "1. Field of the Invention\nThe present disclosure relates to a content-based searching method and system for multimedia objects using high-dimensional feature vector data, which are capable of quickly and exactly processing a query with a 2-level signature.\n2. Description of the Related Art\nDue to developments of computing and media technologies, data are represented in a type of a multimedia including texts, images, audios, and videos. A searching efficiency is a major issue in handling such multimedia information. That is, the major issue is how quickly and exactly a user can search multimedia data containing desired information. A content-based searching method using an extracted high-dimensional feature vector data is generally used for searching data from multimedia objects, such as images, audios, and videos.\nIn the content-based searching method for such multimedia objects, there are two representative query types, that is, a range query and a k-nearest neighbor query.\nThe range query is a query for finding multimedia objects satisfying a similarity with a certain range, and the k-nearest neighbor query is a query for finding the k number of objects with a high similarity. In these query, the calculation of the similarity is generally achieved by a method for calculating Euclidean distance.\nIn order to ensure the fast searching of high-dimensional data, it is important to reduce the calculation of the similarity and the data read. To this end, high-dimensional data indexing schemes have been used, which may be proposed by being largely classified into a tree-based indexing scheme and a filtering-based indexing scheme.\nIn order to efficiently search objects distributed in a data space, the tree-based indexing scheme (e.g., R-Tree, X-Tree, SR-Tree, M-Tree, etc.) uses rectangles or circles indicating a group of neighbor objects as a searching unit. However, since an overlapping area between the rectangles or the circles, which indicate the group of the neighbor objects, is expanded as data dimension increase, there is raised the curse of dimension problem that the searching performance may be lower than that of a sequential searching by exponentially degrading. Therefore, there is a need for a method and system that can solve the curse of dimension problem.\nThe filtering-based indexing scheme (e.g., VA-File, CBF, etc.) is to improve the searching performance for high-dimensional data by using signatures and feature vectors, which is a method that sequentially reads all signature files to primarily filter them and thereafter reads the feature vectors.\nThe filtering-based indexing scheme has a disadvantage that accuracy is lowered when decreasing the bit size of the signature, and that data intended to be read increase when increasing the bit size of the signature."}
{"text": "The present disclosure relates to cable tray systems for use with electrical and other cable and wiring systems.\nCable trays are used to support cables, electrical wires and the like which run through, over, and around buildings. Cable trays may have both straight and curved (horizontal and/or vertical) sections to accommodate the installation requirements of particular settings."}
{"text": "Patterns are recurring design solutions that have been refined over the years by many practitioners to address common design problems. Patterns are usually informally described in the literature for the purpose of education. However, for tools to work with patterns, they need to be formally specified in a machine consumable format. Once formalized, patterns can be either applied or detected in user models.\nPattern detection is one of the fundamental kinds of analysis that can be performed on structured models. Several approaches for pattern detection are available in today's tools based on the way patterns are formalized. In the most rudimentary case, patterns are not formalized and their algorithms are manually coded for every specific pattern. This approach, although flexible, is usually complex to implement, costly to maintain in addition to being error-prone. A better approach is to model an algorithm as a decision tree (Sauve, Alex, “Systems and Computer Engineering,” Masters Thesis, Carleton University, 2005). Although at a higher level of abstraction than code, it still suffers from the same problem, namely the unobvious correspondence between the structure of the pattern and its detection algorithm. The benefits of algorithm flexibility are usually offset by the costs of algorithm validation and maintenance.\nWhen patterns are formalized, detection algorithms tend to be more rigorous. For example, patterns specified with a mathematical notation are usually detected by solvers that employ pattern detection algorithms (see Eden, A. H. et al., “LePUS—A declarative pattern specification language,” Technical Report 326/98, Department of Computer Science, Tel Aviv University, 1998). Although this approach mitigates the need to explicitly think about a detection algorithm, it still suffers from few disadvantages. First, models created in most modeling tools need to be converted to this notation before being analyzed; usually at a performance cost. Another disadvantage is the inherent complexity of the used notation to average modelers.\nAnother approach to pattern formalization is to specify pattern elements as extensions to their domain meta-model elements (see France, R. B. et al., “A UML-Based Pattern Specification Technique, IEEE Transactions on Software Engineering 30(3)193-206, March 2004). Patterns defined in such fashion can be detected by an algorithm that traverses the input model guided by the pattern meta-model. This approach mitigates the conversion problem of the mathematical notation. However, it lacks the ability to configure the detection algorithm. It also lacks a context to the pattern definition that is required to reuse the definition and to represent detected pattern instances. Finally, it forces the definition to have the same complexity as the related part of the domain meta-model complicating pattern detection.\nA new pattern specification formalism called Epattern has been recently defined by the same assignee of the present invention. The new formalism is proposed as an extension of the Ecore meta-model for the purpose of pattern specification. Epattern adds to Ecore new semantics inspired from the composite structure semantics of UML 2.0. As a meta-modeling approach, Epattern focuses mainly on the specification of patterns rather than on their detection algorithms and has simpler semantics than those of the mathematical approaches. What is lacking in Epattern is a configurable pattern detection strategy that integrates with all the semantics of the formalism and allows the pattern author to use some knowledge from the target domain to make it more efficient and scalable.\nPattern specification is a common denominator to most work in applied pattern research. Various approaches have been proposed for pattern specification (Baroni, A. et al., “Design Patterns Formalization”, Ecole Notionale Superieure des Techniques Industrielles, Research Report 03/3/INFO, 2003 and Technical Report 2002). One category of approaches, that Applicant's work also belongs to, uses meta-modeling techniques. The work presented in Guennec, A. et al., “Precise Modeling of Design Patterns,” Proceedings of UML 2000, Vol. 1939 of LNCS, pp. 482-496, Springer Verlag 2000 and Mak, J. “Precise Modeling of Design Patterns in UML,” in Proceedings of the 26th International Conference on Software Engineering, 2004, proposes specifying a pattern as a UML 1.5 meta-collaboration with pattern roles typed with M1 classes stereotyped <<meta>> and named after meta-classes. This obviously prevents writing constraints for such roles as their type information is not available at level M1. Also the binding between a role and an element playing that role is modeled with an explicit dependency relationship. This is in contrast to Applicant's approach which depends on a natural binding between an attribute of a class (the role) and its value (the bound element).\nThe work in Kim, D., “A UML-Based Metamodeling Language to Specify Design Patterns”, in Proceedings of WISME, UML Conference, October 2003, introduces the RBML language, which is used to specify UML patterns as specialized UML meta-models. Pattern roles are specified as sub-classes of their base meta-classes in UML and are related to each other through new meta-associations. One problem with specifying a pattern as a meta-model, rather than a meta-class as in Applicant's approach, is the inability to inherit or compose the pattern which hinders scalability. Another disadvantage is that role binding is done through a generic mapping scheme and is not conveniently an instantiation of the pattern meta-class and an assignment of role values.\nRestated, Applicant finds that use of domain model data (model objects) instead of meta-data (as in the present invention) to describe a pattern limits the expression of the pattern to the semantics of the domain model. This poses a problem if the domain does not have sufficient semantics to completely specify a pattern. Also, this complicates building domain-independent tools that read and process pattern definitions for the purposes of application or detection.\nAnother proposal is found in Maplesden, D. et al., “Design Pattern Modelling and Instantiation using DPML,” in Proceedings of Tools Pacfic 2002, p. 18-21, Sydney, Australia, February 2002, where the DPML language is used to visually specify patterns as a collection of participants, dimensions (multiplicities), relationships and constraints. One drawback is the non-standard notation adopted by the language. Another problem is the restriction of the participants and relationships to predefined types from the UML domain, which limits the scope of the patterns definable by the language. Also, there is no mention of complexity management features.\nAnother approach (Albin-Amiot, H. and Y. G. Gueheneuc, “Metamodeling Design Patterns: Application to Pattern Detection and Code Synthesis”, in Proceedings of the ECOOP 2001 Workshop on Adaptive Object-Models and MetaModeling Techniques, 2001) provides a meta-model to specify patterns. This meta-model is first specialized with pattern domain meta-classes (i.e., meta-classes for every role) before being instantiated to produce an abstract model (pattern specification). Then that model is either instantiated to create a concrete model (pattern instance) or parameterized to use in pattern detection. The provided meta-model contains pattern-domain meta-classes in addition to meta-classes from a target domain (e.g., UML) defined as their subclasses. This need to redefine required meta-classes from the target domain in the pattern meta-model greatly limits the generality and practicality of the approach.\nTo summarize the key differences with Applicant's pattern modeling framework (PMF), most of the above approaches lack the ability to specify patterns for languages other than UML or viewpoints other than the class diagram. They also lack features (e.g., user-defined associations and composition) that help alleviate the complexity of pattern specification. Additionally, some specify M2-level patterns at M1 which deprives them from using free features like pattern constraints and role binding through pattern instantiation. Finally they lack a well-defined process that allows pattern authors the tools to build, refine and simplify patterns in a stepwise manner."}
{"text": "The present invention relates to a shedding apparatus for waste selvage in a loom. More specifically, the invention relates to a shedding apparatus for waste selvage in a loom which has a leno selvage forming apparatus drawing selvage yarns from bobbins and forming a leno selvage.\nIn some fluid jet type looms having a leno selvage forming apparatus for forming a leno selvage, leno selvages are formed in a woven fabric on the weft insertion side of the loom and the other side opposite to the weft insertion side and further waste selvages are formed at positions outward of the respective leno selvages. Generally, in a leno selvage forming apparatus, a planetary gear mechanism including a sun gear and planetary gears is used to cause a pair of yarn guides to revolve so as to draw leno selvage yarns from the respective guides.\nJapanese Unexamined Patent Application Publication No. 2010-100974 discloses a shedding apparatus for catch cord in a loom (shedding apparatus for waste selvage) having a leno selvage forming apparatus in which a planetary gear mechanism is used, as a related art such as shown in FIGS. 9 and 10. As shown in FIGS. 9 and 10, a shedding apparatus for catch cord 100 is provided adjacent to a planetary selvage apparatus 200 (leno selvage forming apparatus) which is provided on the side opposite to the weft insertion side of the loom.\nIn the shedding apparatus for catch cord 100, a beam 204 is provided between opposite side frames (not shown) of the loom and a rail 205 is fixed on an upper surface of the beam 204. A frame 203 of the planetary selvage apparatus 200 is mounted on the upper surface of the rail 205. Specifically, the frame 203 of the planetary selvage apparatus 200 is detachably fixed on the rail 205 with a bolt (not shown) so that the position of the planetary selvage apparatus 200 along the rail 205 can be changed according to a change of the weaving width of a fabric to be woven. A drive shaft 206 extends parallel to the rail 205 on the side thereof that is closer to a let-off side of the loom. The drive shaft 206 is connected to a motor of the loom through a transmission apparatus so that the drive shaft 206 is rotated in conjunction with the spindle of the loom.\nIn the planetary selvage apparatus 200, a support 210 for supporting planetary gears 208 and the like is fixed on a rotating shaft 214 which is rotatably supported by the frame 203. The support 210 is rotated by the drive shaft 206 through a drive timing pulley 209, a driven timing pulley 215, and a timing belt 216. A pair of intermediate gears 218 is engaged with a sun gear 217 at positions angularly spaced from each other at 180 degrees. The number of the teeth of the respective planetary gears 208 is half that of the sun gear 217. Each planetary gear 208 is engaged with its corresponding intermediate gear 218.\nIn the shedding apparatus for catch cord 100, a first shaft 101 and a second shaft 102 are supported by the frame 203 of the planetary selvage apparatus 200. The first shaft 101 supports first swing levers 111, 112 via a first swing lever holder 132, and the second shaft 102 supports second swing levers 113, 114 via a second swing lever holder 133. An eccentric bush 127 whose inner circumferential surface and outer circumferential surface have different axial centers is fitted on the drive shaft 206 at a position that is closer to a warp row than to the drive timing pulley 209. The eccentric bush 127 is fixed to the drive shaft 206 so that the position of the eccentric bush 127 in the axial direction thereof and the position of the eccentric bush 127 in the circumferential direction thereof, i.e., the phase of the eccentric bush 127 may be changed with respect to the drive shaft 206.\nA swing arm 130 is rotatably supported at first end thereof on the eccentric bush 127, and second end of the swing arm 130 and the first swing lever holder 132 are connected to each other via a link pin 134. The first swing lever holder 132 and the second swing lever holder 133 are swingably supported by the first shaft 101 and the second shaft 102, respectively. The first swing lever holder 132 and the second swing lever holder 133 are connected to each other via a connecting link 139.\nWhen the drive shaft 206 is rotated, the first swing lever holder 132 is swung through the eccentric bush 127 and the swing arm 130, thereby causing the second swing lever holder 133 to swing in the reverse direction through the connecting link 139.\nIn Japanese Unexamined Utility Model Application Publication No. 63-106782 discloses a shedding apparatus for waste selvage employing a passive drive system which uses a spring or a cam.\nUnlike the shedding apparatus for waste selvage disclosed in the above Japanese Unexamined Utility Model Application Publication No. 63-106782, the shedding apparatus for waste selvage disclosed in Japanese Unexamined Patent Application Publication No. 2010-100974 employs a positive drive system which does not require a spring or a cam. The shedding apparatus for waste selvage disclosed in Japanese Unexamined Utility Model Application Publication No. 63-106782 that employs a passive drive system, the cam attached to the leno selvage forming apparatus generates a large moment of inertia, resulting in a large increase of the load. The load caused by the spring in the passive drive system is also larger than the load developed in the case of the positive drive system. Therefore, the positive drive system is advantageous over the passive drive system.\nIn the shedding apparatus for wastes selvage disclosed in Japanese Unexamined Patent Application Publication No. 2010-100974, however, the eccentric bush 127 that serves as the drive source for the shedding apparatus for waste selvage (the shedding apparatus for catch cord 100) and the drive timing pulley 209 that transmits drive force to the leno selvage forming apparatus are mounted on the drive shaft 206 that is provided separately from the shaft of the sun gear 217 of the leno selvage forming apparatus (the planetary selvage apparatus 200). Therefore, when changing the weaving width of the loom, fasteners for the drive timing pulley 209 and the eccentric bush 127 need to be loosened to allow the drive timing pulley 209 and the eccentric bush 127 to be relocated and then be re-tightened after the relocation is over. Furthermore, the timings need to be set again for each apparatus. Thus, the work associated with changing of the weaving width is troublesome and time consuming.\nThe present invention, which has been made in view of the circumstances above, is directed to providing a shedding apparatus for waste selvage in a loom which reduces the load applied to the drive for the shedding apparatus for waste selvage and requires less time for the work associated with changing the weaving width of the loom."}
{"text": "The present invention relates to a blocking device for a wedge-type breechblock in a gun, and more particularly to a blocking device of the type which is arranged parallel to the bore axis of the gun and which includes a detector. A recess is provided in a base ring at the base of the gun barrel to accommodate an obturating ring, and the detector engages the recess under the force of a spring if the obturating ring is absent from the recess.\nSuch a blocking device is disclosed in U.S. Pat. No. 3,420,139. In such a blocking device, if the obturating ring is missing, the detector drops into the recess for the obturating ring and interrupts the locking process of the breechblock.\nThe kinetic energy of the closing breechblock, which must be decelerated at once, may cause damage to the sealing system of the base ring, particularly at the contact faces of the obturating ring. Repair of such damages consumes time and causes a considerable amount of installation and adjustment expenditures."}
{"text": "The X-ray transmission technology is commonly used in inspection of cargos in a container in places such as airports, customs etc. The X-ray transmission technology utilizes a linear detector array to obtain a two-dimensional transmission image of the container. For example, X-rays are collimated to form a fan-shaped ray beam, which is transmitted through an inspected object for detection. When the object is stationary, the detector receives only a one-dimensional projection of the object on a ray layer. When the object moves in a direction perpendicular to the ray layer, all parts of the object sequentially pass through the ray beam, and the detector obtains a two-dimensional projection image of the whole object by line scanning. However, in the inspection manner described above, projections of the object in the container in a depth direction overlap together, which results in a loss of depth information. Especially for materials with light atomic numbers such as explosives and dangerous articles etc., grayscale attenuation after X-ray transmission is not obvious enough, which seriously affects recognition and identification of features such as a shape, a position etc. of such an object.\nCommon measures which are currently available for acquiring a three-dimensional transmission image are CT tomography and stereo matching technology etc. The CT tomography technology is complicated and expensive, and has a long imaging time, which limits its application in industrial real-time detection. In a process of stereo matching for an complex object, grayscale reconstruction and boundary processing are complex, the calculation complexity is also very large, and therefore the stereo matching is difficult to be applied in practice. Therefore, it is necessary to study other technologies to acquire three-dimensional imaging information of the object in the container."}
{"text": "The present invention relates to apparatus for storing and dispensing particulate material and in particular to a toner or developer cartridge for electrostatographic printing machines.\nIn an electrostatographic reproducing apparatus commonly in use today, a photoconductive insulating member is typically charged to a uniform potential and thereafter exposed to a light image of an original document to be reproduced. The exposure discharges the photoconductive insulating surface in exposed or background areas and creates an electrostatic latent image on the member which corresponds to the image areas contained within the usual document. Subsequently, the electrostatic latent image on the photoconductive insulating surface is made visible by developing the image with developing powder referred to in the art as toner. Most development systems employ a developer material which comprises both charged carrier particles and charged toner particles which triboelectrically adhere to the carrier particles. During development the toner particles are attracted from the carrier particles by the charge pattern of the image areas in the photoconductive insulating area to form a powder image on the photoconductive area. This image may subsequently be transferred to a support surface such as copy paper to which it may be permanently affixed by heating or by the application of pressure. Following transfer of the toner image to a support surface, the photoconductive insulating member is cleaned of any residual toner that may remain thereon in preparation for the next imaging cycle.\nAs the toner particles are depleted from the developer material it is necessary to dispense additional toner particles into the developer mixture. Further, the carrier has a limited life due to a variety of problems occurring with continued use. For example, the carrier may become impacted with toner thereby altering its triboelectric properties or if the carrier is coated with the material to enhance triboelectric properties the coating may deteriorate with time. U.S. Pat. No. 4,614,165 to Folkins et al. describes a developing process involving the addition of both toner particles and carrier particles to the developer in the developer housing to insure that the usable life of the developer material in the chamber at any point in time is at least equal to the life of the photographic printing machine. To accommodate the addition of further carrier and toner material, waste or spent developer is removed from the developer housing when it seems to exceeds a predetermined quantity. Thus, in most commercial embodiments the concentration of toner particles within the developer mixture is maintained substantially constant. To achieve this particulate material containers for toner and/or developer which discharge the toner and/or developer into the development system have been used. In replenishing the toner particles in an electrostatographic printing machine, it is important to minimize any spillage which may result in contamination of otherwise uncontaminated areas. Further, the toner particles being very finely ground may become airborne carrying this contamination to other areas but immediately adjacent the development system. In addition, the spilled toner particles have a tendency to cling to an operator's hands or to the surrounding environment. Accordingly, there is a need to provide a toner and/or developer storing and dispensing device which may be used to periodically replenish the toner and/or developer in an automatic printing machine and which can avoid the above difficulties thereby overcoming a dirty and messy operation.\nWhile the desire to minimize contamination by the toner or developer and provide for \"white glove\" machine operator involvement during the replenishment activity of the toner or developer is of paramount importance it is also noted that the toner and/or the developer container must be capable of containing the toner and/or developer during periods of storage and/or shipment under conditions which may go from one extreme to another by way of temperature, pressure and/or impact. A common approach to sealing a toner container or cartridge has been to seal the dispensing opening with a breathable plastic film of polyvinyl chloride, for example, which has a pressure and heat actuated adhesive on one side which upon the application of heat and pressure forms a very good seal with the dispensing opening. When necessary for use, the dispensing opening can be opened and particulate materials emptied into the developer housing. A principle difficulty, however, associated with this type of sealing arrangement is the preciseness required for pressure, temperature and time in sealing the toner cartridge. For example, if a hot iron is left on to long, it may burn or melt through the plastic film. Accordingly, it has been found that the plastic film is generally too time temperature dependent. In addition, the adhesive is frequently applied unevenly and with the uneven application of the adhesive in trying to manually peel or remove the film from the opening the force that may have to applied may be uneven and an ineffective seal may be formed in areas where too little adhesive is applied. Further, once this plastic film seal is broken and upon exhaustion of the supply of toner and/or developer in the cartridge so that it must be removed and replenished with a full cartridge, there is really effectively no seal to protect the user from toner spillage and contamination upon loading the cartridge from the machine."}
{"text": "1. Field of the Invention\nThe present invention relates in general to doors for a motor vehicle, and more particularly, to automotive doors of a type which is equipped with a power window.\n2. Description of the Prior Art\nHitherto, various automotive doors of the above-mentioned type have been proposed and put into practical use. However, due to their inherent constructions, some of them have failed to make a satisfied operation of the power window mounted therein. In fact, after long use, some of them tend to have such a problem that the upward and downward movement of the window pane can not be made smoothly due to unavoidable deformation of the door. The deformation of the door tends to cause a play of the window pane particularly at the time when the window pane is halfly or incompletely opened. Such play produces a noise during cruising of the vehicle. Furthermore, some of the automotive doors have failed to have a sufficient space for accommodating the power window. In this case, the work for mounting the power window to the door is difficult.\nFurthermore, hitherto, inspection of the power window has been carried out after the same is properly mounted to the door. However, such inspection is troublesome. In fact, when a defect is found in the power window, it becomes necessary to dismantle the same from the door for repairing the same. Thereafter, the power window should be mounted to the door again. This is very troublesome and time consuming."}
{"text": "Mobile devices, such as, cellular phones and personal digital assistants (PDAs), are often configured with short range wireless transmitters to enable wireless communication over a network. The signals transmitted by the wireless transmitters may be detected by a base station when the device is within the proximity of the base station. A base station is a radio receiver/transmitter that serves as the hub of a local wireless network and may also be the gateway to a wired network."}
{"text": "The present invention relates to new derivatives of N-(2,2-dichlorovinyl)salicylamide and use thereof for controlling bacteria, fungi and algae which propagate on industrial raw materials, industrial products and in water in circulating water systems.\nThe bacteria and fungi on industrial raw materials and products degrade the quality of industrial raw materials and products and cause various other troubles. The damage is very serious. The algae which propagate on the bottom of ships cause the slowdown of speed of ships.\nBacteria, fungi and algae in water in circulating water systems produce mats of slime which restrict the flow of water and reduce the efficiency of heat exchange. They constitute a considerable problem.\nMany compounds having bactericidal, fungicidal and algicidal activity are used for prevention of the damage caused by the bacteria, fungi or algae to industrial raw materials and products and the impediment caused by slime in circulating water systems. These compounds include, for example, organic mercury compouds, organic tin compounds, phenols having substituted halogen atoms, organic sulfur compounds, formalin, cresol and quaternary ammonium salts. These compounds, however, have various defects, for example, strong toxicity against humans and live-stock, pollution of environment, irritative or offensive odor, bad influence upon industrial raw materials and products and weak bactericidal, fungicidal activity."}
{"text": "1. Field of the Invention\nThis invention relates generally to automatic transmissions for trucks, and more specifically to automatic truck transmissions fitted with bypass valves for reducing dead time or delay in transmitting fluid control pressure signals when gears are changed manually.\n2. Brief Description of the Prior Art\nThe Detroit Diesel Allision Division of General Motors Corporation is the largest producer of automatic transmissions for truck and other heavy vehicles. An important product has been the MT 650 transmission with five forward speeds. This transmission was introduced in model year 1974. The construction, operation, and repair of this transmission are described in detail in Service Manual, Allison Transmissions, Automatic Models MT 640, 643, MT 650, 653 (Detroit Diesel Allision, Indianapolis, 1983), incorporated herein by reference. The MT 653 transmission is a modified version of the 650 transmission.\nIn the early production models of the MT 650 and 653 transmissions an extended time delay of from about five to forty seconds of clutch-up or engagement for reverse gear operation occured in some transmission applications. The delay was a result of the priority design of the hydraulic control system and occured when a shift was made from low gear to reverse gear. In the MT 650 and 653 transmissions, low gear, the lowest forward gear, is a very low gear which must be engaged manually, and first gear is the lowest forward gear which is shifted automatically. There are always two clutches engaged in the drive train for each gear. In first gear these are the first clutch and the forward clutch; in reverse gear these are the first clutch and the fourth clutch. In neutral, there is only one clutch applied, the first clutch. Thus, the shift from low gear to reverse gear involve releasing the forward and low clutches and applying the first and fourth clutches.\nIn an attempt to remedy the delay, the transmissions were modified by installing a bypass fitted with a check valve in between the reverse pressure line (fourth clutch circuit) and the first clutch apply pressure line (first clutch circuit). The check valve included a plastic ball valve member and was positioned to seal when the pressure in the first clutch apply pressure line was greater than the pressure in the reverse pressure line, and to open when the pressure in the reverse pressure line was greater than the pressure in the first clutch apply pressure line.\nThe plastic ball of the check valve would tend to dissolve i the transmission fluid. After a period of about three to five years, in particular in transmissions equipped with power-take off (\"PTO\") fittings for supplying hydraulic pressure to hydraulic cylinders used for sundry purposes in certain types of heavy vehicles, the ball would fail to seal when the transmission fluid pressure was raised in the first clutch pressure line.\nIn these transmissions, a forward clutch and the low clutch are simultaneously applied in operation of the first forward gear, and the first clutch and a fourth clutch are simultaneously applied in operation of reverse gear. However, if, when first gear is engaged and the second clutch pressure line is then pressurized, the check valve fails to seal, the reverse pressure line is also pressured, and the forward, first and fourth clutches are all applied simultaneously, resulting in severe damage to the transmission. This requires a time-consuming and expensive rebuilding of the failed transmission. If the ball were to fail when the transmission was in neutral, the first and fourth clutches would be unintentionally applied, and the vehicle would begin to move backwards, a potentially very dangerous situation.\nA number of different plastic materials for the ball, including nylon, were apparently substituted by the transmission manufacturer in attempts to solve this substantial problem, all without ultimate success.\nSince the time that installation of the plastic check valve in new model transmissions was initiated by the manufacturer, the manufacturer also altered the transmissions by replacing the input oil pressure pump with a \"wide gear\" version so as to increase the rate at which input oil pressure increases when the engine is accelerated Previously, the drive and driven gears of the oil pump had a width specified to be from 0.6835 inch to 0.6855 inch while the \"wide gear\" pump which replaced the earlier model has drive and driven gears which have a width specified to be from about 0.8180 inch thick to 0.8205 inch thick Two different driven gear widths and four different drive gear widths have been provided by the manufacturer for replacement, as slightly wider gears must be fit when a transmission is rebuilt and the oil pump cavity has become worn through circulation of oil suspending abrasive particles. The width of the drive gear is specified to be within plus or minus two and one half ten thousands of an inch and the with of the driven gear is specified to be within plus or minus five ten thousands of an inch.\nThe transmission manufacturer provided instructions for eliminating the delay in units in field service in \"S.I.L. 10-TR-80,\" Feb. 20, 1980, as well as in Instruction Sheet No. 230, used in connection with Service Kit Part Number 23041150, including a template and replacement separator plate for use in these transmissions. S.I.L. 10-TR-80 and Intruction Sheet No. 230 are each incorporated herein by reference.\nHeavy vehicles equipped with Allison MT 650 transmissions are widely used, and the rebuilding required after about three to five years service is extremely costly for owners of fleets of such vehicles, including utility companies. The costs include not only the direct cost of rebuilding the transmissions, but also the cost of maintaining larger fleets than would otherwise be necessary, the costs to customers for whom service is delayed because of transmission failure, and the like. Despite apparently repeated efforts to remedy the problem, the manufacturer has apparently been unsuccessful. There is a significant need for an improved transmission which simultaneously overcomes the problem of delay encountered when shifting from first to reverse gear, and the premature failure encountered after only about three to five years of service of transmissions having plastic ball check valves in between the first clutch apply pressure line and the reverse pressure line."}
{"text": "The present invention relates in general to coupled line devices and is concerned, more particularly, with quadrature hybrids and couplers including directional couplers constructed in accordance with the principals of the present invention and having in particular improved directivity and power handling capabilities. In accordance with the present invention the device has preferred high directivity, low power consumption and high thermal transfer.\nU.S. Pat. No. 3,358,248 shows a previous version of a coupled line device including a pair of insulated inner conductors and a common outer conductor with the inner conductors spaced a distance corresponding substantially to a quarter wavelength at the center operating frequency. If the parallel transmission line center conductors could be imbedded in a uniform dielectric material, whether air or some other material, the even and odd mode propagation velocities will be equal. However, in practice this does not occur and in the device described in U.S. Pat. No. 3,358,248, the conductors are not imbedded in a material having a uniform dielectric. For example, the coating on the wires typically has a relatively high dielectric constant which may be in the order of 2.7. In the previous construction which employed a Teflon support bead, the Teflon has a dielectric constant in the order of 2.0. This inconsistency in dielectric constant slows the odd mode propagation in comparison to the even mode. Also, the wires have a twist which can make the electrical length larger for the odd mode than for the even mode. Futhermore, there tend to be air voids between the wires and the surrounding Teflon. These air voids, depending upon location, may either increase or decrease the odd mode velocity compared to the even mode velocity.\nAccordingly, it is an object of the present invention to provide an improved miniaturized coupled line device in which the propagations delay of both the even and odd modes are substantially equalized.\nAnother object of the present invention is to provide an improved coupler in accordance with the preceding object and which is characterized by high directivity over a wide frequency range.\nAnother object of the present invention is to provide a coupled line device that may be constructed as a TEM directional coupler and which is characterized by relatively high power handling capabilities.\nAnother object of the present invention is to provide a coupled line device that may be constructed as a 3 dB hybrid in which is characterized by relatively high powered handling capabilities.\nStill another object of the present invention is to provide a high performance miniaturized coupled line device that is relatively easy and inexpensive to fabricate.\nA further object of the present invention is to provide an improved fabrication technique for a microwave coupler that enables all air voids to be filled thus improving operation and also preventing moisture entry into the device.\nAnother object of the present invention is to provide an improved microwave coupler device that is constructed to furthermore provide improved center conductor cooling resulting from the filling of air voids."}
{"text": "It is known to provide chairs, especially stackable chairs as described in the Mauser office brochure entitled \"Mahrzweckstuhle Serie 570\" (Multipurpose Chair Series 570) of Mauser Office GmbH, Korbach, Germany, in which the chairs are described for auditoriums and the like which can be aligned in rows and connected together by coupling means between adjacent chairs. These chairs can be fitted with so-called lecture desks, i.e. a table carried by a mounting or holding mechanism enabling the table to be releasably attached to the chair so that the table is disposed just forwardly of the user and, ordinarily, on one side of the chair.\nThe lecture desk can be swung from a more or less horizontal position of use in which it is locked above the seating surface and forwardly of the user, into a substantially folded position alongside the seating service. The holding mechanism can, if desired, be attached to an armrest of the chair.\nFor assembling the chairs into rows for use in an auditorium or the like, the stackable chairs are removed from the stack, positioned side by side and connected together by appropriate coupling means. When the chairs are taken up from the rows, they can be stacked in a nested state so that space is conserved.\nThe coupling systems used for interconnecting the chairs in rows, are available in various configurations. For example, German patent 39 33 817 discloses a coupling system in which a stirrup is provided on the underside of the seating surface of one chair and hooks into a stirrup of the underside of an adjoining chair when these stirrups are extended toward one another. The free ends of the stirrups are formed with mutually-engaging hooks, A problem with this arrangement is that whether the chair is provided with or without an armrest, the lecture desk must be configured differently and different configurations of the couplings may have to be used.\nEarlier desk systems have not proved to be effective for chairs which could be provided with or without armrests, are adapted to be coupled together or to be free-standing, or are especially configured for assembly in rows."}
{"text": "This invention relates to automated packaging machines for small round produce, such as blueberries and cherry tomatoes, and in particular a dispensing mechanism for an automated conveyer driven produce packaging system.\nThin plastic containers have become increasingly popular as containers for small produce, such as blueberries, cherry tomatoes, and other small fruit. These containers include a xe2x80x9cclam shellxe2x80x9d lid, which is integrally attached to one edge of the container and thus may be moved from an open position to a closed position in which the corners of the lid engage the corners of the receptacle to thereby close the container until it is opened by the consumer.\nAutomated systems have been developed for packaging small produce in clam shell containers. Conventional packaging systems include a variety of automated conveyers and mechanisms, which perform various packaging functions. Generally conveyers are used to transport produce and containers through the various automated work stations, where an automated mechanism performs some packaging operation, such as metering produce into the container or closing the containers after it is filled with produce. While automated systems reduce the time and man power required to get produce to market, cost, complexity and lack of flexibility present several practical and operational draw backs for small produce growers and packagers.\nConventional automated packaging systems use large gravity feed dispensing apparatus. The dispensing apparatus generally include a large hopper and an adjustable gate, which meters the flow of produce from the hopper into the containers, as the containers are carried along a conveyer. The size of the gate opening can be adjusted to increase or decrease the volume of produce flowing into the containers. Generally the gates are located at the bottom of the hoppers, which are suspended directly over the conveyers such that the produce falls directly downward into the containers. The mouth of the gate is elevated above the containers only a short distance to prevent the produce from being damaged by the fall.\nConventional gravity feed dispensing apparatus have had several draw backs. Conventional dispensing mechanisms often damage the produce as it is metered into the containers. Since the metering gates of overhead dispensing hoppers are located directly over the containers and conveyers, the weight of the bulk produce inside the hopper exerts significant pressure on the produce at the mouth of the gate. The pressure on the produce can shear and bruise the produce as it is force through the gate opening. The location of the hoppers directly above the conveyer mechanism, can also contribute to produce damage. If the gate opening is not adjusted to generate the appropriate flow rate, the containers may be over filled. Since the gate is elevated above the containers only a short distance, overfilling the container will result in the excess produce being scrapped off and out of the container by the sides of the gate, as the container is transported along the conveyer. Over filling not only results in shear damage and loss of produce, but may also prevent the container from being closed.\nConventional dispensing mechanisms fail to consistently meter the same volume of produce into each successive container on the conveyer. The location of metering gates in overhead dispensing hoppers creates a xe2x80x9cplowing effect,xe2x80x9d which prevents conventional dispensing mechanisms from metering a consistent volume of produce. The size of individual pieces of produce (berries and the like) varies greatly. As containers pass under the dispensing gate of an overhead hopper, large individual pieces of produce, caught in the gate and not sheared off by the gate of an overhead hopper, block or xe2x80x9cplowxe2x80x9d smaller pieces of produce from entering the containers. As a result, the volume of produce metered into each successive container varies greatly, with some containers being under filled and others being over filled. In addition, locating the hopper and gate directly over the conveyer also makes regulating the volume and flow rate of the produce through the gate more difficult. The pressure on the produce at the mouth of the gate exerted by the weight of the bulk produce fluctuates with the volume of produce inside the hopper. While the hopper is a gravity feed system, the fluctuation of the pressure exerted on the produce can influence the flow rate of the produce into the containers. Often, the volume of produce inside the hopper is not consistently maintained. As a result, the gate opening must be constantly adjusted to regulate the flow rate of the produce into the containers. While this problem can be reduced by using conveyers to deposit a steady stream of bulk produce into the dispensing hopper, additional conveyers add to the cost and complexity of the packaging system.\nRecently, dispensing hoppers for produce packaging machines have been developed to address some of these problems. For example, one packaging system mechanism manufactured by AandB Blueberry Packaging of Hartford, Mich., includes a side located dispensing hopper, which meters produce through a gate in the container at an angle, instead of vertically from an over head hopper. Other dispensing hoppers have included internal baffles, which divert the flow of produce and reduce the pressure on the produce at the mouth of the metering gate. Nevertheless, these apparatus have not adequately addresses the problems of gravity feed dispensing apparatus.\nThe dispensing mechanism of this invention is used for dispensing produce into clam shell type containers as part of an automated packaging system. The dispensing mechanism is mounted to a packaging machine that includes a conveyer mechanism for transporting the containers past the dispensing mechanism. The dispensing mechanism is mounted to a subframe, which extends over the conveyer track of a conventional packaging machine. The dispensing mechanism of this invention includes a side mounted hopper, a guide shoot, and a pair of elongated produce stops. The interior of the hopper is divided by internal partitions into a large bin and a dispensing queue. The body of the hopper has an inclined floor, and converging side and end walls that funnel produce toward two vertically oriented metering gates. One gate meters produce from the large bin into the queue and the second gate meters produce from the queue into the containers. The guide shoot is pivotally connected to the hopper beneath the second gate. The guide shoot directs the produce from the hopper into the containers, and also is used to adjust the speed of produce from the hopper into the containers. The produce stops are located in front of the shoot and run directly over and substantially parallel to the conveyer track to prevent the loss of produce and direct the produce into the center of the container base. The stops can be adjusted horizontally and vertically to insure that the containers are filled evenly and without loss of produce.\nThe dispensing mechanism combines several unique features to address the draw backs of the conventional produce dispensing apparatus. The use of a divided hopper interior and two vertically oriented metering gates, along with the side mounted location of the hopper, allows produce to be dispensed into containers with less damage to the produce and with greater control of the flow rate. The angular attitude of the hopper and the shoot can be independently adjusted to vary the speed of the produce moving through the hopper and over the shoot, respectively. The weight of the bulk produce is exerted directly onto the metering gates. Produce flow through the gate openings at an angle to the direction of gravitational force; consequently, the force exerted on the produce passing through the mouth of the gate is significantly less and the flow rate is less effected by the volume of produce inside the hopper. Dividing the interior of the hopper into a large accumulating bin and a smaller dispensing queue divides the total weight of the produce in the hopper into two areas. The first gate meters produce from the bin into the queue to maintain a relatively constant volume of produce inside the queue. Maintaining a constant volume of produce inside the queue, allows the shoot gate to consistently meter constant flow rate from the hopper, which is unaffected by the volume of the produce inside the bin. In addition, the smaller volume of produce inside the queue, compared to the larger volume in the bin, exerts less pressure on the produce, which has less influence on the flow rate from the hopper into the containers. Consequently, the flow rate through the gate from the queue can be consistently maintained and controlled.\nAccordingly, an advantage of this invention is to provide for a dispensing mechanism for an automated conveyer driven packaging system, which dispenses produce at a consistently controlled flow rate to evenly fill each successive container with a uniform volume of produce.\nAnother advantage of this invention is that the dispensing mechanism reduce damage to the produce.\nAnother advantage of this invention is that the dispensing mechanism reduces the effect of fluctuations in the volume of bulk produce in the hopper on the flow rate of produce into the containers.\nAnother advantage of this invention is that each functional component can be readily adjusted to control the volume, speed and flow rate of the produce.\nOther advantages will become apparent upon a reading of the following description."}
{"text": "When it has been desired to cut a series of holes in a picture mat, this has generally required a laborious and time consuming manual effort. Such efforts are prone to miscalculations or mismarkings which can ruin the job.\nIn order to cut ovals in picture mats, there has been disclosed in the prior art an Oval-Master apparatus manufactured by C&H Manfacturing of Jackson, Miss. This device was adjustable to cut ovals of various sizes and of ovalities including plane circles as well. Using a suitable movable guide on a base, the picture mat was suitably located beneath a center of the cutting apparatus. The cutting apparatus was mounted to a platform which was lowerable into position to engage the mat during cutting. While this apparatus was advantageous for cutting ovals and circles in mats, in order to cut a plurality of such ovals and circles in a mat, the mat had to be continually repositioned at precise locations determined by x and y scales on the base."}
{"text": "A terahertz wave is an electromagnetic wave having a frequency of approximately 0.01 THz to 100 THz corresponding to an intermediate range between light waves and radio waves, and has an intermediate property between light waves and radio waves. As an application of such a terahertz wave, a technology for acquiring information on a measuring object by measuring a temporal waveform of an electric field amplitude of a terahertz wave which is transmitted through or is reflected by the measuring object has been studied.\nA technology for measuring information on a measuring object by use of a terahertz wave is generally as follows. That is, pulsed light output from a light source (for example, a femtosecond laser light source) is split into two components to be pulsed excitation light and pulsed detection light by a branching part. The pulsed excitation light among those is input to a terahertz wave generating element (for example, a nonlinear optical crystal or a photoconductive antenna element), and a pulsed terahertz wave is generated from the terahertz wave generating element thereby. This generated terahertz wave is transmitted through or reflected by a measuring object part as a measurement light, to acquire information on the measuring object (for example, an absorption coefficient, a refractive index), and thereafter, the terahertz wave is made incident on a terahertz wave detecting element (for example, an electro-optic crystal or a photoconductive antenna element) in substantially the same timing as that of the pulsed detection light.\nIn the terahertz wave detecting element to which the terahertz wave and the pulsed detection light are input, a correlation between both light components is detected. For example, in the case in which an electro-optic crystal is used as the terahertz wave detecting element, the terahertz wave and the pulsed detection light are coupled by a coupling part, to be made incident on the electro-optic crystal, and birefringence is induced in accordance with propagation of the terahertz wave in the electro-optic crystal, and a polarization state of the pulsed detection light is changed by the birefringence. A change in the polarization state of the pulsed detection light in the electro-optic crystal is detected, and as a result, an electric field amplitude of the terahertz wave is detected, which enables to acquire information on the measuring object.\nTerahertz wave generating apparatus of various configurations for generating terahertz waves have been known. One of those terahertz wave generating apparatus is configured to make pulsed excitation light incident into a nonlinear optical crystal to generate a terahertz wave from the nonlinear optical crystal. In such a conventional terahertz wave generating apparatus, a type of nonlinear optical crystal suitable for terahertz wave generation is determined according to a central wavelength of the pulsed excitation light. For example, in the case where pulsed excitation light with a central wavelength of 800 nm is used, a ZnTe crystal is broadly used as a nonlinear optical crystal on the basis of a phase matching condition for the pulsed excitation light and the terahertz wave in the nonlinear optical crystal. When pulsed excitation light with a central wavelength of 800 nm is incident into a ZnTe crystal, a pulsed terahertz wave is generated coaxially with the pulsed excitation light.\nIn recent years, it has become apparent that, even in a combination of pulsed excitation light of a central wavelength and a nonlinear optical crystal in which a phase matching condition is not usually satisfied, the pulsed excitation light whose pulse front is tilted is made incident into the nonlinear optical crystal, to enable to satisfy the phase matching condition (refer to Non Patent Documents 1 to 3)."}
{"text": "1. Field of the Invention\nThe present invention relates to the flat panel TV (especially the LCD TV) display system, and more specifically to a video system for improving the response time and reducing the motion-blur of the panel display.\n2. Description of Prior Art\nThe development of flat panel displays has been accelerated. In the current display market, flat panel displays are gradually replacing the traditional CRT displays. However, some of the flat panel displays, for example, TFT-LCD panel displays, are generally slower in the response time compared to the CRT displays. This disadvantage hampers the development of the LCD TV, for which the display quality (especially motion blur) is highly concerned.\nTo improve the quality of motion picture on the flat panel (especially TFT-LCDs), various driving techniques have been introduced. Among them, the over driving method is well known as the effective scheme for improving the response time. The over driving method makes the response time fast by applying an over driving voltage to a motion pixel. In order to apply the over driving method to the flat panel, the motion pixels are detected by comparing the pixel values of current and previous frames, and determining the value to be added to the original value. The over driving level is decided through a decision model. Such a model heavily depends on the characteristics of the flat panel. For TFT-LCDs, the decision model is normally implemented by a look-up table built in a over drive circuit chip. The over drive decision requires one-frame delay mechanism for getting the pixel value of the previous frame. A compression and decompression mechanism is used for reducing the bandwidth demands of the frame buffer access for implementing the one-frame delay. The over driving method is proposed in some prior arts to combine with the timing controller (TCON) in the flat panel display for improving the video quality of motion picture.\nTo further improve the response time of a flat panel display, the present invention proposes an extra video processing for de-motion-blurring. A video system in accordance with the present invention combines the color TV decoding, deinterlacing, scaling, and de-motion-blur process and many other video enhancement features to optimize the overall system in cost."}
{"text": "Commercial processing of thermoplastic polyers, such as extrusion and molding, is well-known in the prior art and may be accomplished by a diverse variety of techniques. Generally, thermoplastic polymers have softening temperatures of about 130.degree. C. and higher and, therefore, require processing at temperatures significantly higher, such as 200.degree. C. and higher.\nIn the past, these 200.degree. C. and higher processing temperatures were of little concern because the necessary heat could be generated from a variety of heat sources and capital equipment. However, with the current worldwide energy shortage and the increasing concern related to the worker's environment, there is a need for a commercial process which permits high-softening temperature thermoplastic polymers to be processed at temperatures lower than those currently used. Such a lower processing temperature represents a significant reduction in energy used, a significant savings in energy costs and an improvement in the worker's environment since less heat is likely to be dissipated by the processing equipment.\nFurthermore, it is desirable to prepare thermoplastic polymers which contain additives, such as chemical blowing agents, thermostabilizers, hydrolytic stabilizers, mold release agents, pigments, etc. which are molded or extruded at a later date or at a different site. When incorporating additives in such thermoplastic polymers, it is desirable to use as low a temperature as possible in order to minimize the heat histories of such additives and thermoplastic polymers in order to avoid deterioration of the physical, mechanical or chemical properties of the thermoplastic polymer or additives.\nA severe problem currently exists in the preparation of a masterbatch or concentrate of a thermoplastic polymer intimately blended with a chemical blowing agent. It is often easier first to prepare such masterbatches or concentrates and then subsequently add them to a larger quantity of a compatible thermoplastic polymer. However, due to the processing temperature necessary to soften many thermoplastic polymers, a chemical blowing agent cannot be intimately blended with the polymer to form a concentrate without activating the chemical blowing agent. This problem, therefore, often requires the physical blending of such a chemical blowing agent with a thermoplastic polymer at the actual molding site making this molding step more expensive, cumbersome and difficult due to the powder form of many chemical blowing agents.\nThus, a real need exists for a process which permits high softening temperature, thermoplastic polymers to be processed at temperatures lower than those currently used.\nIn accordance with the present invention, a method is provided which permits high softening temperature, thermoplastic polymers to be processed at temperatures lower than those currently used. Novel chemical blowing agent concentrates are further provided."}
{"text": "Scene execution in a home automation system relates to preconfiguring a scene, or a set of actions, to be taken by one or more specified devices in the system when one or more events occur. For example, a vacation scene can include instructing lighting devices in the system to be turned on at a first set of predetermined times and to be turned off at a second set of predetermined times.\nBypassing a zone in a home automation system relates to deactivating one or more zones in a region monitored by the system before arming the system. Once the system is armed, the system will not monitor bypassed zones.\nExecuting scenes and bypassing zones are known in the art. For example, in known home automation systems, scenes are executed and zones are bypassed based on predetermined schedules or manual user input. However, these known systems fail to account for the location of users in a region monitored by the home automation system.\nIn view of the above, there is a continuing, ongoing need for improved systems and methods."}
{"text": "In recent years, the weight reductions and thinning of display apparatuses like flat panel displays (FPDs: liquid crystal display apparatuses, organic EL display apparatuses, and the like) have been progressing with the development of visual communication technologies. Glass substrates have been conventionally used as substrates for use in the display apparatuses in many cases. The glass substrates are each excellent in transparency, solvent resistance, gas barrier properties, and heat resistance. However, when one attempts to achieve the thinning of a glass material of which any such glass substrate is formed, the following problem arises. While the weight reduction is achieved and the glass substrate shows excellent flexibility, its impact resistance is insufficient, and hence the glass substrate becomes difficult to handle.\nIn order that the handleability of thin glass substrates may be improved, substrates in each of which a resin layer is formed on a glass surface have been disclosed (see, for example, Patent documents 1, 2, and 3). Such substrate is produced by attaching a resin film to the glass surface or coating the glass surface with a thermoplastic resin. The substrate thus produced can obtain such a certain effect as to have flexibility.\nHowever, when such substrate is used in a display device or the like, the following problem arises. The substrate is poor in durability against various solvents (such as acetone) to be used upon production of the display device."}
{"text": "1. Field of the Invention\nThis invention relates to an eye fundus photographing apparatus for photographing the fundus of an eye to be examined.\n2. Related Background Art\nUsually, when photographing the fundus of an eye to be examined having a small pupil diameter, it is more advantageous to photograph at a narrow field angle of a high magnification than to photograph at a wide field angle of a low magnification. However, when photographing at a narrow angle of field, as compared with the photographing at a wide angle of field, the quantity of photographing light is reduced and therefore, in Japanese Patent Publication No. 63-36251, a baffle comprising a circular light intercepting plate is installed in an illuminating optical system with a view to prevent the harmful effects of reflected light from the front face of the cornea or the rear face of the crystalline lens, and as the field angle is made narrower, the diameter of the baffle is made smaller so as to increase the quantity of photographing light.\nHowever, in the above-described example of the prior art, when photographing an eye to be examined having a small pupil diameter, it is required to make the baffle and the diameter of an illuminating field stop smaller than in the ordinary case and therefore, a special condition must be set and a new operation must be added."}
{"text": "1. Field of the Invention\nThe invention relates to a torque detection device that detects, for example, a steering torque of a steering wheel and an electric power steering system that includes the torque detection device.\n2. Discussion of Background\nA torque sensor described in US 2006/0137474 A is formed of a torsion bar, a ring-shaped magnet, a pair of annular magnetic yokes, a pair of magnetic flux concentration rings, a magnetic sensor, and the like. The torsion bar couples an input shaft to an output shaft such that the input shaft and the output shaft are coaxial with each other. The magnet is connected to an end portion of the input shaft. The magnetic yokes are connected to an end portion of the output shaft. The magnetic flux concentration rings are arranged in proximity to the outer peripheries of the magnetic yokes. Each of the magnetic yokes has lugs all around. The number of the lugs is equal to the number of N-poles and S-poles of the magnet. The magnetic flux concentration rings concentrate magnetic fluxes induced from the magnet by the magnetic yokes. The magnetic sensor is placed between the magnetic flux concentration rings, and detects the magnetic flux density generated due to the magnetic fluxes concentrated by the magnetic flux concentration rings.\nWhen torque (steering torque) is input into a portion between the input shaft and the output shaft through a steering operation of a steering wheel, the torsion bar is twisted, and the relative position between the magnet and the magnetic yokes in the circumferential direction changes. The torque sensor detects the steering torque input into the portion between the input shaft and the output shaft on the basis of the magnetic flux density that changes in the magnetic yokes in accordance with the change in the relative position.\nIn the torque sensor described in US 2006/0137474 A, a portion of each magnetic yoke, at which each lug is provided, has an L-shape such that, when viewed from the circumferential direction of the magnetic yoke, the portion extends radially inward of the magnetic yoke, bends at a substantially right angle and then extends in the axial direction of the magnetic yoke (see FIG. 5 of US 2006/0137474 A). The distal end portions of the lugs of the respective magnetic yokes are arranged alternately in the circumferential direction. In the torque sensor, the lugs through which magnetic fluxes flow are desired to be elongated by extending the lugs in the axial direction in order to improve torque detection capability and torque detection accuracy. However, if the size of each lug in the axial direction (hereinafter, in the specification, referred to as “effective length” where appropriate) is elongated excessively, the distal end of the lug excessively approaches the opposed magnetic yoke. Thus, the magnetic yokes may interfere with each other portions close to each other. In addition, the magnetic flux is likely to unexpectedly leak from the magnetic yoke at that portion of each lug. Therefore, it is difficult for the magnetic sensor to highly accurately detect the magnetic flux density in the magnetic yokes (the density of magnetic fluxes concentrated by the magnetic flux concentration rings). That is, it is difficult for the torque sensor to highly accurately detect a torque.\nIn addition, it is always desired to reduce the size of the torque sensor."}
{"text": "The advantages of sustained release pharmaceutical compositions are well-known and their use is becoming more and more widespread. The methods used in this field of pharmaceutical industry are continuously developping. However the requirements raised against sustained release compositions are considerably higher and more severe than those raised against conventional pharmaceutical compositions. For this reason there is not yet known any process which can be applied for every active ingredient and meets all requirements. The preparation of sustained release pharmaceutical compositions from active ingredients having a relatively high therapeutical dose is particularly encountered with difficulties because the solid oral pharmaceutical compositions are of a limited size (the weight of a single dosage unit form can not be generally more than 0.8-1.0 g) and for this reason the known methods wherein the protracted effect was achieved with the aid of a relatively higher amount of auxiliary agents (the amount of the said auxiliary agent was approximately the same as that of the active ingredient) were applicable not at all or but to a very limited extent.\nSustained release pharmaceutical compositions corresponding to a relatively high single dosage unit of 100-500 mg can be prepared first of all in the form of matrix tablets. This is due to the fact that pharmaceutical compositions which have the lowest specific volume and may be hence most easily swallowed can be manufactured by pressing.\nAccording to U.S. Pat. No. 2,895,881 waxes or fatty substances or hydrogenated castor oil can be used as release retarding agent in the preparation of matrix tablets. For this purpose glycerol monostearate (U.S. Pat. No. 2,993,839), a mixture of stearic acid and castor oil (U.S. Pat. No. 2,736,628), waxes, a mixture of waxes and water insoluble substances and mixture of waxes and hydrophylic polymers (U.S. Pat. Nos. 4,132,753, 3,402,240, 3,459,850 and 3,487,138) can also be applied. According to U.S. Pat. Nos. 2,987,445 and 3,317,394 water insoluble polymers, polyethylene and polymethyl-methacrylate and polivinyl chloride can be used as release retarding agent.\nIt is also known that a mixture of waxes and water insoluble polymers (U.S. Pat. No. 3,965,256), mixtures of polymers which readily swell in water and form mucus (U.S. Pat. No. 3,065,143), polymers forming complex with each other in water (carbopol; an acrylic acid polymerisate) and polyvinyl pyrrolidone (U.S. Pat. No. 3,458,622) and a mixture of carbopol and polyoxyethylene glycole (U.S. Pat. No. 3,634,584) can also be used as release retarding agent. In British Pat. No. 935,672 the use of protein derivatives while in U.S. Pat. No. 3,062,720 the use of water insoluble substances (e.g. talc, calcium sulfate, calcium hydrogen phosphate etc.) is reported. According to U.S. Pat. No. 3,905,508 a mixture of talc, ethyl cellulose and a metal stearate is used.\nThe matrix tablets disclosed in the said references comprise 10-70% of active ingredient. Sustained release tablets having a higher active ingredient content than the above value can only be prepared from active ingredients being poorly soluble in aqueous medium (to less than 1%).\nAccording to the cited patent specifications matrix tablets are prepared by means of conventional granulating methods (by wet or dry granulation) or in the case of fatty substances by melting the said materials and encorporating the active ingredient into the liquid matrix material."}
{"text": "1. Field of the Invention\nThe present invention relates to detachable cushions and more particularly pertains to a portable seat cushion designed to be attached directly to a user's clothing.\n2. Description of the Prior Art\nThe use of detachable cushions is known in the prior art. However, these cushions are generally designed to be attached to existing seat structures and not to the body or clothing of a particular user. A good typical example of a prior art detachable cushion is to be found in U.S. Pat. No. 4,212,496 which issued to A. Kirkham, Jr. on Jul. 15, 1980. This patent discloses a detachable cushion for the seat of ski lift chairs and is usable for the purpose of providing some measured degree of comfort and dryness to a chair user. However, when a ski lift chair is not being utilized, it is continually exposed to ambient weather conditions which effects both its comfort and dryness. As such, there appears to exist a need for some type of cushion which could be utilized on a ski lift chair or under similar conditions wherein such cushion could be maintained in a dry condition until its use is necessitated. In this respect, the present invention substantially fulfills this need."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of computer systems and, in particular, to a method and apparatus for stress testing of a circuit board assembly.\n2. Background Information\nManufacturers have long understood that manufacturing processes are not perfect. Manufacturers of printed circuit board assemblies (PCBs), for example, understand that the manufacturing processes used to manufacture PCBs are unable to produce a population of product which performs exactly the same. Therefore, manufacturers determine performance specification limits for each critical operating parameter of a product, generally drawn around an optimum performance level for each critical parameter. For example, the operating frequency of a clock source may be specified at 66 MHz, with a \"specification window\" of plus or minus 2 MHz. Thus, a clock source that performs within the range of 64-68 MHz will be identified as \"acceptable\" (or \"passing\"), while those not performing within this range will be identified as \"unacceptable\" (or \"failing\"). The purpose of the specification window is to ensure the proper performance of the PCB assembly when integrated into larger systems, and under the operating conditions of the end-user environment.\nNonetheless, despite this \"weeding out\" of unacceptable product through the parametric testing which occurs during the manufacturing process, there remains a percentage of the \"acceptable\" population of product that fails to meet performance specifications when subjected to the actual operating conditions of the end-user environment. One explanation for the variability in performance of the printed circuit board assembly may be attributed to the variability in the components of the PCB. Thus, if a product just barely falls within the specification window for a particular operating parameter, i.e. the product is marginally acceptable, it may well be that under certain operating conditions the product may fail to meet the operating specification.\nConsequently, one of the issues faced by manufacturers is how to identify and thus contain the marginally performing population of products produced by the manufacturing process that will fail under the conditions of the end-user operating environment. One prior art solution to this problem is to simply tighten the specification window. This is a very costly solution, however, in that not all of the products tested to perform near the edge of the specification window will actually result in a field failure.\nAnother prior art solution to the problem is subjecting the population of product, or a sample thereof, to environmental stress testing. Environmental stress testing, as its name implies, subjects the circuit board assembly to a number of environmental stresses such as temperature (both hot and cold), humidity, and shock/vibration testing.\nEnvironmental stress testing, however, has the limitation that it employs a mechanical delivery mechanism. That is, environmental testing applies the stress to the electronic device externally, through the packaging of the devices. The packaging of electronic devices has a property commonly referred to as a \"thermal coefficient\" which describes the rate at which heat is transferred through the packaging material. Without going into unnecessary details, it is sufficient to say that environmental stress testing incurs a time delay, necessary to allow the inner circuitry of the electronic devices to settle into a steady state operating temperature, before the parametric testing may be completed. Accordingly, environmental testing is also considered a costly solution, both in terms of the time required to complete the testing sequence, as well as the capital equipment associated with conducting such tests (i.e., temperature chambers, humidity chambers, shock/vibration chambers). Therefore, despite manufacturers effort to identify and contain marginally performing printed circuit board assemblies, the prior art solutions have not proven to be effective.\nThus a need exists for a more cost effective method and apparatus for stress testing of a circuit board assembly that is unencumbered by the deficiencies and limitations associated with the prior art."}
{"text": "Technical Field\nThis application generally relates to data storage, and more particularly to techniques used in connection with providing data storage hints.\nDescription of Related Art\nComputer systems may include different resources used by one or more host processors. Resources and host processors in a computer system may be interconnected by one or more communication connections. These resources may include, for example, data storage devices such as those included in the data storage systems manufactured by EMC Corporation. These data storage systems may be coupled to one or more host processors and provide storage services to each host processor. Multiple data storage systems from one or more different vendors may be connected and may provide common data storage for one or more host processors in a computer system.\nA host processor may perform a variety of data processing tasks and operations using the data storage system. For example, a host processor may perform basic system I/O operations in connection with data requests, such as data read and write operations.\nHost processor systems may store and retrieve data using a storage device containing a plurality of host interface units, disk drives, and disk interface units. Such storage devices are provided, for example, by EMC Corporation of Hopkinton, Mass. The host systems access the storage device through a plurality of channels provided therewith. Host systems provide data and access control information through the channels to the storage device and storage device provides data to the host systems also through the channels. The host systems do not address the disk drives of the storage device directly, but rather, access what appears to the host systems as a plurality of logical disk units, logical devices, or logical volumes (LVs). The logical disk units may or may not correspond to the actual disk drives. Allowing multiple host systems to access the single storage device unit allows the host systems to share data stored therein.\nIn connection with data storage, a variety of different technologies may be used. Data may be stored, for example, on different types of disk devices and/or flash memory devices. The data storage environment may define multiple storage tiers in which each tier includes physical devices or drives of varying technologies, performance characteristics, and the like. The physical devices of a data storage system, such as a data storage array, may be used to store data for multiple applications."}
{"text": "Light Amplification by Stimulated Emission of Radiation (laser) provides for a spatially coherent, spatially narrow, spectrally narrow (monochromatic), low divergence beam of electromagnetic radiation. Lasers can be found in a number of modern-day applications, such as cutting tools, surgery, optical storage devices, and optical communications. A laser beam is produced by reflecting light back and forth using reflectors through an optical cavity containing a gain medium having an energy population inversion (i.e. more electrons in a relatively high energy state versus at a low energy state). The reflectors used for the reflection are typically flat, spherical, or parabolic. Additionally, once the laser beam is emitted through a partially transmitting reflector, the beam can be manipulated by optical elements such as reflectors, lenses, polarizers, optical fiber, and diffraction gratings. For some applications, such optical elements may need to be very precisely manufactured and may have a resulting high cost. Lenses, in particular may need to have a very precise shape and surface finish for some applications. The material properties of a lens, such as the index of refraction, vary with the wavelength of the laser beam. The uses of lenses to converge a beam of light may lead to various optical aberrations, including spherical and chromatic aberrations."}
{"text": "The present invention relates to a data structure suitable for recording/playing back video data sent in digital broadcast or the like or stream data sent with a packet structure, a method of recording stream data using this data structure, and a method of playing back stream data recorded with this data structure.\nIn recent years, TV broadcast has come into the era of digital broadcast. Accordingly, an apparatus for saving digital data of digital TV broadcast as they are irrespective of their contents, i.e., a so-called streamer, has been demanded.\nThe current digital TV broadcast uses an MPEG transport stream. In the future, an MPEG transport stream is projected to be used as a standard one in the field of digital broadcast using moving picture.\nIn such digital broadcast, the contents (mainly, video information) to be broadcasted are time-divided into groups of data each having a predetermined size (e.g., 188 bytes) called transport packets, and broadcast data is sent in units of transport packets.\nAs a streamer for recording digital broadcast data, a home digital VCR such as D-VHS (digital VHS) or the like is currently commercially available. A streamer using D-VHS directly records a broadcasted bitstream on a tape. For this reason, a plurality of programs are multiplexed and recorded on a video tape.\nUpon playback, all data are output from the VCR to a set-top box (digital TV reception apparatus; to be abbreviated as an STB hereinafter) when they are played back either from the beginning or the middle of the tape. In this STB, a desired program is selected from the output data by user operation or the like. The selected program information is transferred from the STB to a digital TV receiver, and is played back (playback of video+audio and the like).\nSince this D-VHS streamer uses a tape as a recording medium, it cannot attain quick random access, and it is difficult to quickly jump to a desired position of a required program so as to play it back.\nAs a promising candidate that can combat such shortcoming (difficulty of random access) of the tape, a streamer that uses a large-size disc medium such as a DVD-RAM or the like has been proposed. In this case, management data must be inevitably recorded together with broadcast data in consideration of random access, special playback, and the like.\nProblem\nDigital TV broadcast adopts an MPEG transport stream as a transmission method of stream data containing video information, and video information is sent while being combined in units of application packets each having, e.g., 188 bytes. By contrast, when a recording medium such as a DVD-RAM disc or the like that belongs to a DVD family is used, recording must be done in units of sectors, the minimum recording unit of which is 2,048 bytes. However, the following problems remain unsolved.\n(1) A method of efficiently recording information of 188-byte application packets in units of 2,048-byte sectors is not settled yet.\n(2) A method of playing back stream data recorded on an information storage medium while holding the reception timing of digital TV broadcast is not settled yet.\nFurthermore, there is also a need for recording digital data which is transferred with a packet structure in a local area network (LAN) or a digital telephone line such as ISDN or the like using a streamer.\nHowever, the following problem also remains unsolved.\n(3) No versatile recording method that can record digital data other than digital TV is available.\nObject\nIt is an object of the present invention to provide a data structure (recording format) on an information storage medium, which can efficiently record stream data on an information storage medium, and can play back stream data from an information storage medium while holding transfer timings between application packets received in digital TV broadcast.\nIt is another object of the present invention to provide a method of recording stream data using the data structure.\nIt is still another object of the present invention to provide a method of playing back stream data recorded with the data structure.\nIn order to achieve the above object, in a data structure according to the present invention, a set of one or more stream objects (SOB) that represent playback data for a recorded bitstream constitutes stream data, the stream object (SOB) is formed of one or more stream packs (S_PCK), and the stream pack (S_PCK) is formed of a pack header and stream packet (S_PKT). The pack header includes predetermined time information (SCR), and the stream packet (S_PKT) includes one or more application packets (AP_PKT) assigned with predetermined time stamps (ATS). The application packet (AP_PKT) is incoming (to a streamer) during recording of the stream object (SOB) is time-stamped (modified time stamp in FIG. 14; ATS in FIG. 15) by a local reference clock (440 in FIG. 9) (in the streamer) corresponding to the predetermined time information (SCR).\nIn order to achieve the other object, in a recording method according to the present invention, a set of one or more stream objects (SOB) that represent playback data for a recorded bitstream constitutes stream data, the stream object (SOB) is formed of one or more stream packs (S_PCK), and the stream pack (S_PCK) is formed of a pack header and stream packet (S_PKT). The pack header includes predetermined time information (SCR), and the stream packet (S_PKT) includes one or more application packets (AP_PKT) assigned with predetermined time stamps (ATS). The application packet (AP_PKT) is incoming (to a streamer) upon recording the stream object (SOB) on an information medium (201) is time-stamped (modified time stamp in S4 in FIG. 10; time stamp in ST106, ST212, and ST312 in FIGS. 21 to 23) by a local reference clock (440 in FIG. 9) (in the streamer) corresponding to the predetermined time information (SCR).\nIn order to achieve the still other object, a playback method according to the present invention uses an information medium (201) on which stream objects (SOB) are recorded in a format in which a set of one or more stream objects (SOB) that represent playback data for a recorded bitstream constitutes stream data, the stream object (SOB) is formed of one or more stream packs (S_PCK), the stream pack (S_PCK) is formed of a pack header and stream packet (S_PKT), the pack header includes predetermined time information (SCR), the stream packet (S_PKT) includes one or more application packets (AP_PKT) assigned with predetermined time stamps (ATS), the incoming application packet (AP_PKT) (to a streamer) is time-stamped by a local reference clock (440 in FIG. 9) (in the streamer) corresponding to the predetermined time information (SCR). When recorded information is played back from this medium, a playback reference clock is set (S37 in FIG. 11) on the basis of the local reference clock (SCR 303 in FIG. 5, SCR base in FIG. 15) played back from the information medium (201), and contents of the bitstream are played back from the information medium (201) on the basis of the set playback reference clock (SCR).\nAdditional objects and advantages of the invention will be set forth in the description which follows, and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out hereinafter."}
{"text": "This invention relates to the transmission and reception of digital data over radio links, and more particularly, to the transmission and reception of 4-level and 8-level signaling symbols according to prevailing conditions so as to achieve a greater data rate when conditions allow.\nIn a typical cellular telephone system, speech compression is used in order to accommodate as many mobile users as possible within the spectrum allocated to any particular operator. Speech compression involves a tradeoff between the total number of users that can be served by the system, and the speech quality delivered by the system. Where less speech compression is used, a better speech quality is achieved; however, less users can be served by the system. Where a greater amount of speech compression is used, the total number of users that can be served by a system is increased; however, the speech quality is degraded. Similarly, there is a tradeoff between the quality of a data communication service and the number of users that can be accommodated by the service. The quality of a data service may be measured in terms of the time required to transfer a data file. Thus, using a higher bit rate may succeed in completing data transfer faster, thereby leading to increased quality of data communications; however, other users are then prevented from using the spectrum at the same time. Additionally, the number of users that can be accommodated in the spectrum may be increased; however, the quality of data communications is degraded. Thus, an improved system would be one which is able to use less speech compression or use a higher bit rate to transmit data, without reducing the number of users that can be accommodated by the system.\nFurther, in a cellular telephone system, a transmitted signal may travel from a transmitter to a receiver over multiple paths, for example, a direct path and one or more reflected paths. In digital cellular phone systems, an equalizer and an error correction decoder are used to process the received signal which may have traveled over multiple paths. In situations where the received signal contains two modulation types, for example, 4-level modulation and 8-level modulation, two equalizers and two error correction decoders are required to process the received signal, where the first equalizer and first error correction decoder are suitable for processing 4-level modulation and the second equalizer and second error correction decoder are suitable for processing 8-level modulation. This increases the hardware cost of a cellular station since multiple equalizers and error correction decoders are required. Furthermore, in a field where products are being manufactured increasingly smaller for convenience, the user of multiple equalizers and error correction decoders require greater printed circuit board (PCB) real estate, resulting in a physically larger product than cellular stations using only one equalizer and error correction decoder.\nThe present invention is directed to overcoming one or more of the problems discussed above in a novel and simple manner.\nIn accordance with the invention, there is disclosed a method and system for the transmission of both 4-level and 8-level signalling symbols using an 8-level modulator. There is further disclosed a method and system for demodulating and error correction decoding a signal containing both 4-level and 8-level modulation.\nIt is an object of the present invention to provide a method for modulating 4-level signaling symbols using an 8-level modulator, where the method comprises receiving 4-level signaling symbols to be modulated, where each 4-level signaling symbol contains 2-bits of information. The method further includes expanding the 4-level signaling symbols into eight-level signaling symbols, where each 8-level signaling symbol contains 3-bits of information. The method further includes modulating the 8-level signaling symbols.\nIn a preferred embodiment, the step of expanding the 4-level signaling symbols into 8-level signaling symbols includes manufacturing a third bit for each 4-level signaling symbol and appending the third bit to the 2-bits of information. In this preferred embodiment, the third bit is the exclusive-OR of the 2-bits. In another preferred embodiment, the method includes performing a progressive 45xc2x0 phase shift between a modulation of consecutive symbols. In a further preferred embodiment, modulating and expanding 4-levels signaling symbols with the 8-level modulator causes the 8-level modulator to produce 4-level modulation. In a yet further embodiment of the invention, the 4-level signaling symbols are a subset of the 8-level signaling symbols. In another preferred embodiment of the invention, the method includes the step of sending an indication signal indicating which of expanded 4-level signaling symbols and 8-level signaling symbols are being transmitted.\nIt is another object of the present invention to provide a method for demodulation of a received signal containing at least one of 4-level modulation and 8-level modulation where the method includes demodulating the received signal under a first assumption that the received signal contains 4-level modulation and determining a value of a first quality factor by performing error correction decoding on symbols modulated under the first assumption. The method further includes demodulating the received signal under a second assumption that the received signal contains 8-level modulation and determining a value of a second quality factor by performing error correction decoding on symbols modulated under the second assumption. Additionally, the method determines a received modulation present in the received signal using the first and the second quality factors, and demodulates the received signal as the received modulation.\nIn a preferred embodiment, the step of determining from the first and second quality factors may be performed before the demodulation under the first and the second assumptions is completed. In a further preferred embodiment, the first and second quality factors are determined using metrics generated by an error correction decoder.\nIt is another object of the invention to provide a method for demodulating a received signal containing at least one of 4-level modulation and 8-level modulation, where the method determines whether 4-level modulation or 8-level modulation is present in the received signal. The method continues by operating a demodulator in a constrained 4-level demodulation mode by restricting possible levels for the received signal to four levels, determining which level out of the four levels is being received, and producing 4-level signaling symbol decisions having 2-bits each if it is determined that 4-level modulation is present. If however, it is determined that 8-level modulation is present in the received signal, the method continues by operating the demodulator in an unconstrained 8-level demodulation mode by determining which level out of 8 possible levels is being received, and producing 8-level signaling symbol decisions having 3-bits each.\nIn a preferred embodiment, the step of determining whether 4-level modulation or 8-level modulation is present includes the step of receiving an indication signal indicating whether 4-level modulation or 8-level modulation is present in the received signal. In a further preferred embodiment, the indication signal is transmitted using at least one of a SYNCWORD, a CDVCC, and a FACCH message. In a further preferred embodiment, the step of determining whether 4-level modulation or 8-level modulation is present includes the step of using a predetermined transmission format to determine whether a 4-level signaling symbol or an 8-level signaling symbol is present in the received signal.\nIt is another object of the invention to provide a method for decoding a received signal containing at least one of 4-level signaling symbols and 8-level signaling symbols, the method comprises receiving a signal, demodulating the received signal as 8-level modulation, and determining whether 4-level signaling symbols or 8-level signaling symbols are present in the received signal. If it is determined that 4-level signaling symbols are present in the received signal, the method operates an error correction decoder in a constrained 4-level decoding mode and produces 4-level signaling symbol decisions having 2-bits each. If it is determined that 8-level signaling symbols are present in the received signal, the method operates the error correction decoder in an unconstrained 8-level mode and produces 8-level signaling symbol decisions having 3-bits each.\nIn a preferred embodiment, the step of determining whether 4-level signaling symbols or 8-level signaling symbols are present further includes the step of receiving an indication signal indicating whether 4-level signaling symbols or 8-level signaling symbols are present in the received signal. In a further preferred embodiment, the step of determining whether 4-level signaling symbols or 8-level signaling symbols are present further includes the step of using a predetermined transmission format to determine whether 4-level signaling symbols or 8-level signaling symbols are present in the received signal. In another preferred embodiment, the step of decoding the received signal in the unconstrained 8-level mode includes removing a third bit from the 3-bits of the 8-level signaling symbol decision.\nIt is another object of the invention to provide a method for decoding a received signal containing at least one of 4-level signaling symbols and 8-level signaling symbols, where the method comprises receiving a signal, demodulating the received signal as 8-level modulation, and performing error correction decoding under a first assumption that 4-level signaling symbols are present in the received signal and determining a first quality factor. The method continues by performing error correction decoding under a second assumption that 8-level signaling symbols are present in the received signal and determining a second quality factor. The method then performs the step of determining a received modulation in the received signal using the first and the second quality factors, and decodes the received signal as the received modulation.\nIn a preferred embodiment, the step of determining from the first and second quality factors may be performed before the decoding under the first and the second assumptions is completed.\nIt is another object of the invention to provide a system for modulating 4-level signaling symbols using an 8-level signaling symbol modulator, where the system comprises a processor for processing information into a plurality of 4-level signaling symbols to be modulated. The system also includes an expander coupled to the processor for expanding the plurality of 4-level signaling symbols each containing 2-bits of information to expanded 8-level signaling symbols each containing 3-bits of information. The system further includes an 8-level modulator coupled to the expander for modulating the expanded 8-level signaling symbols.\nIn a preferred embodiment, the system applies a progressive 45xc2x0 phase shift between the modulation of consecutive symbols. In a further preferred embodiment of the system, the 4-level signaling symbols are a subset of the expanded 8-level signaling symbols. In another preferred embodiment of the invention, the 8-level modulator produces 4-level modulation when modulating the expanded 8-level signaling symbols.\nIt is yet another object of the invention to provide a system for demodulating a received signal containing at least one of 4-level modulation and 8-level modulation, where the system comprises a receiver for receiving the received signal, and a demodulator coupled to the receiver. The demodulator operates under at least one of a first assumption that the received signal contains 4-level modulation, and a second assumption that the received signal contains 8-level modulation. If the demodulator is operating under the first assumption that the signal contains 4-level modulation, the demodulator operates as a constrained demodulator by restricting the possible levels of the received signal to a set of four levels, determining which level out of the four levels is received, and producing 4-level signaling symbol decisions having 2-bits each. If the demodulator is operating under the second assumption that the received signal contains 8-level modulation the demodulator operates as an unconstrained demodulator by restricting possible levels of the received signal to a set of eight levels, determining which level out of the eight levels is received, and produces 8-level signaling symbol decisions having 3-bits each.\nIn a preferred embodiment, the system may further include a first error correction decoder coupled to the demodulator and a second error correction decoder coupled to the demodulator, where the first error correction decoder decodes the 4-level signaling symbol decisions when the demodulator is operating under the first assumption, and the second error correction decoder decodes the 8-level signaling symbol decisions when the demodulator is operating under the second assumption. In a further preferred embodiment, a controller is coupled to the first and the second error correction decoders and to the demodulator where the first error correction decoder determines a first quality factor, the second error correction decoder determines a second quality factor, and the controller uses the first and the second quality factors to determine whether 4-level modulation or 8-level modulation is present in the received signal. If the controller determines that 4-level modulation is present, the controller directs the demodulator to operate as the-constrained demodulator, and if the controller determines that 8-level modulation is present, the controller directs the demodulator to operate as the unconstrained demodulator. In-yet another further preferred embodiment, a controller is coupled to the demodulator where the controller determines whether 4-level modulation or 8-level modulation is present in the received signal by using at least one of an indication signal within the received signal and a predetermined transmission format.\nIt is another object of the invention to provide a system for decoding a signal containing at least one of 4-level modulation and 8-level modulation, where the system comprises a receiver for receiving a signal, and a demodulator coupled to the receiver for demodulating the signal as 8-level modulation. The system also includes an error correction decoder coupled to the demodulator for operating as at least one of a constrained 4-level mode for performing error correction under a first assumption that 4-level signaling symbols are present in the signal, and an unconstrained 8-level mode for performing error correction decoding under a second assumption that 8-level signaling symbols are present in the signal.\nIn a preferred embodiment, the system comprises a controller coupled to the error correction decoder wherein the error correction decoder determines a first quality factor when operating under the first assumption, and a second quality factor when operating under the second assumption, and the controller determines a received modulation from the first and the second quality factors. The controller then sends a control signal to the error correction decoder to decode the signal as the received modulation. In another preferred embodiment of the invention, the error correction decoder removes a third bit from the 3-bits created by the demodulator when performing error correction decoding under the second assumption. In a further preferred embodiment, the system includes a controller coupled to the error correction decoder wherein the controller determines whether 4-level or 8-level signaling symbols are present in the demodulated signal, and the controller directs the error correction decoder to operate in the constrained mode if it is determined that 4-level signaling symbols are present, and directs the error correction decoder to operate in the unconstrained mode if it is determined that 8-level signaling symbols are present. In a further preferred embodiment, the controller uses at least one of an indication signal within the received signal and a predetermined transmission format to determine whether 4-level or 8-level signaling symbols are present."}
{"text": "The present invention relates to a biogas plant and a servicing device for a biogas plant.\nVarious biogas plants and servicing devices for biogas plants have become known in the prior art. Biogas plant are as a rule provided with one or more fermenter tanks for fermenting the fermenting substrates and as a rule with an agitator provided for stirring and at least partially homogenising the fermenting substrates in the fermenter tank.\nMost of these agitators have agitator blades with which to mix and agitate the fermenting substrates. In dependence on the substrates to be fermented and further ambient conditions and on the materials employed the agitators require inspection and if necessary maintenance at specified intervals. Some or all of the agitator blades may be exchanged for example.\nFor maintenance on an agitator the agitator which is disposed on a mounting or holding rod or the like is traversed from the lower fermentation section in which the fermenting substrates are located, to the upper gas section. For gaining access to the agitator, for example a biogas plant equipped with a flexible top sheeting may be opened and the top sheeting removed entirely to thus gain easy access to the agitator. This method, however, involves the drawback of considerable work for removing the top sheeting. Also, the biogas collecting in the gas section above the fermentation section escapes as the top sheeting is removed.\nFor biogas plants having a concrete top, DE 199 51 959 A1 has made known a rectangular maintenance shaft inserted in the concrete top allowing maintenance of the agitator without significantly compromising the gas production. In the case of fermenters having a high solids content the maintenance shaft is bipartite with the top part terminating approximately 10 cm above the normal liquid level. In the case of maintenance a lower shaft portion is inserted into the top shaft portion for immersing in the fermenter mass. For sealing the top shaft portion is provided with a circumferential sealing channel with which when assembled a sealing collar of the lower shaft portion engages. These two-part maintenance shafts involve the drawback that each case of maintenance requires newly inserting the lower shaft portion and sealing against the top shaft portion and removing again after maintenance. Also, the fermenter filling volume must be kept within narrow limits for achieving reliable sealing. Therefore this document proposes for fermenters having a low solids content a one-piece maintenance shaft reaching at least 10 cm beneath the normal liquid filling level. For balancing the gas a shut-off pipe connection to the gas headspace of the fermenter is provided. In the case of maintenance the work involved is reduced compared to a two-piece maintenance shaft. The drawback of this is, however, that these maintenance shafts are suitable for low solids contents only. Also, deposits may accumulate on the immersed shaft. Another drawback is a permanent increase in the agitating resistance and thus in energy consumption.\nEP 1 717 305 B1 has made known a biogas plant servicing shaft for a fermenter of a biogas plant with a top sheeting with the servicing shaft being provided as a gas-tight, dome-shaped cover across a servicing opening in the fermenter tank cover. The agitating unit held vertically in the fermenter tank can be moved upwardly from the lower fermentation section into the servicing shaft interior. After opening a servicing flap at the servicing shaft the agitator device can be serviced in the interior of the dome-type servicing shaft and one or all of the agitator blades may be exchanged as needed.\nA biogas plant servicing shaft according to EP 1 717 305 B1 operates reliably, allowing servicing the agitator device while involving little work wherein maintenance does not require removal of the entire top sheeting. However, as the servicing flap is opened, the biogas which has collected in the upper gas section of the fermenter tank and which is under overpressure, escapes while air and thus a considerable amount of oxygen enters the fermenter tank requiring removal prior to restarting operation."}
{"text": "(1) Field of the Invention\nThe present invention relates to a liquid crystal display device, and in particular, to the form of a sealing portion through which a TFT substrate and a facing substrate are pasted together in a liquid crystal display panel.\n(2) Related Art Statement\nIn liquid crystal display devices, the space between the TFT substrate on which pixel electrodes and thin film transistors (TFT's) are formed in a matrix and the facing substrate on which color filters are formed is filled in with liquid crystal, so that molecules of the liquid crystal can be controlled through an electrical field, and thus, images can be formed. In general, a sealing material is applied around the periphery of the facing substrate using a sealing material applying apparatus, referred to as dispenser, and the TFT substrate and the facing substrate are pasted together.\nIn the case where the space between the substrates is filled in with liquid crystal in accordance with a suction method, the TFT substrate and the facing substrate are pasted together with a sealing material in between and the sealing material is cured, and after that, the inside space is vacuumed and filled with liquid crystal. Meanwhile, in the case where the space between the substrates is filled in with liquid crystal in accordance with a dropping method, liquid crystal is dropped on the TFT substrate before the TFT substrate and the facing substrate are pasted together, and then the facing substrate is pasted to the TFT substrate. After that, the sealing material is cured.\nIn the case where a sealing material is provided using a dispenser, the sealing material is applied in annular form around the periphery of the facing substrate in such a manner that the starting point of application using a dispenser overlaps with the finishing point, for example. “Patent Document 1” and “Patent Document 2” describe a technology for forming the finishing point or the starting point of application as described above at a distance from the display region toward the outside so that the overlapping portion does not affect the display region.    (Patent document 1) Japanese Unexamined Patent Publication 2003-241204    (Patent Document 2) Japanese Unexamined Patent Publication 2002-122870"}
{"text": "1. Field of the Invention\nThe present invention relates to charge transport compositions. The invention further relates to photoactive electronic devices in which there is at least one active layer comprising such charge transport compositions.\n2. Background\nIn organic photoactive electronic devices, such as light-emitting diodes (“OLED”), that make up OLED displays, the organic active layer is sandwiched between two electrical contact layers in an OLED display. In an OLED the organic photoactive layer emits light through the light-transmitting electrical contact layer upon application of a voltage across the electrical contact layers.\nIt is well known to use organic electroluminescent compounds as the active component in light-emitting diodes. Simple organic molecules, conjugated polymers, and organometallic complexes have been used.\nDevices which use photoactive materials, frequently include one or more charge transport layers, which are positioned between the photoactive (e.g., light-emitting) layer and one of the contact layers. A hole transport layer may be positioned between the photoactive layer and the hole-injecting contact layer, also called the anode. An electron transport layer may be positioned between the photoactive layer and the electron-injecting contact layer, also called the cathode.\nThere is a continuing need for charge transport materials and anti-quenching materials."}
{"text": "The December 2012 shooting at former Sandy Hook Elementary School was one of a number of prominent shooting incidents to occur in schools across the United States over the last 20 years. Doors to classrooms are often not equipped with a locking device that may be locked from the inside. Therefore, students located inside a classroom are unable to quickly prevent the doors from opening if an armed assailant wants to enter a classroom."}
{"text": "The invention relates to planarization of a semiconductor wafer by polishing the wafer with a polishing pad and a polishing fluid.\nAn integrated circuit comprises many semiconductor transistors that are electrically connected to one another by circuit interconnects constructed as tiny metal lines. The metal lines transmit electrical signals that operate the transistors. The transistors and circuit interconnects are of miniature size, and are constructed according to semiconductor manufacturing processes that involve formation of the transistors and circuit interconnects by depositing or growing different semiconductor materials on a wafer of silicon, i.e., a semiconductor wafer.\nA semiconductor wafer is adapted for manufacture thereon of many circuit interconnects arranged according to a pattern. Circuit interconnects are constructed on a semiconductor wafer, by depositing a dielectric layer on an underlying silicon wafer, followed by recessing tiny trenches in the surface of the dielectric layer, the trenches being arranged in a pattern, followed by applying a known barrier film on the dielectric layer, followed by depositing a layer of metal onto the underlying barrier film and dielectric layer to assure that the trenches are filled with metal. A process known as chemical mechanical planarization, CMP, is performed with a polishing pad and a polishing fluid to remove the layer of metal and the barrier film from the underlying dielectric layer, and to provide a polished surface that is smooth and planar. Thus, the process of CMP is often described as a method of polishing a semiconductor substrate, or as a method of planarizing a semiconductor substrate.\nAn aim of the process of CMP is to polish the metal in the trenches to a height that is substantially the same height as that of the smooth, planar polished surface. The polished metal in the trenches provides the circuit interconnects. The dielectric layer with the circuit interconnects provides a substrate on which successive dielectric layers are constructed. The successive dielectric layers have additional circuit interconnects to provide successive layers of circuit interconnects for integrated circuits that are constructed in large numbers on the patterned wafer.\nPlanarizing a semiconductor substrate by CMP is performed with a polishing pad and a polishing fluid, which together remove material from the surface of the substrate. The polishing pad rubs against the surface to remove material by abrasion. Simultaneously, the polishing fluid is dispensed at an interface of the polishing pad and the substrate. The polishing fluid chemically reacts with the surface, and dissolves material that results from chemical reaction of the surface with the polishing fluid. Thus, material is removed from the surface by both mechanical abrasion and chemical reaction.\nThe surface has a surface topography with high topography features, i.e., high spots, and low topography features, i.e., low spots. The wafer is planarized, by CMP polishing of the surface with a polishing pad and a polishing fluid, which removes material from the topography features, until the topography features become substantially the same height. However, the polishing fluid reacts chemically with both the high spots and the low spots, and removes material from both of them. Thus, polishing the high spots and the low spots to the same height is counteracted by having the polishing fluid react chemically with the low spots.\nThe metal circuit interconnects in the trenches have corresponding low spots in the metal layer prior to a polishing operation. When such low spots on the circuit interconnects are polished by CMP, too much of the metal can be removed, as indicated by the presence of excessive concavity in the polished surfaces of the circuit interconnects. Dishing is a term that refers to the concavity in the circuit interconnects, as caused by CMP. Dishing is a structural defect of the circuit interconnects that contributes to unwanted attenuation of the signals transmitted by the circuit interconnects. Accordingly, a CMP process must be performed with care to minimize dishing.\nAttempts have been undertaken to hasten the process of planarization by CMP, which would lower manufacturing costs. However, increasing the speed of a CMP process to rapidly remove material from the surface of a semiconductor wafer, causes rapid dishing, and causes an uneven polished surface, which is indicative of deficient planarization. Thus, a need exists for an invention that minimizes dishing, while minimizing the time duration for planarizing a surface by polishing with a polishing pad and a polishing fluid.\nThe invention minimizes the time duration for planarizing a surface on a semiconductor wafer by polishing the surface with a polishing pad and a polishing fluid.\nThe invention adapts a semiconductor wafer for removal of high topography features at a faster removal rate than for low topography features. The invention further adapts a semiconductor wafer to minimize dishing of circuit interconnects, while the high topography features are removed by polishing with a polishing pad and a polishing fluid.\nThe invention adapts a semiconductor wafer for minimized removal rate of the low topography features, while the high topography features are removed by polishing with a polishing pad and a polishing fluid.\nFurther, the invention adapts a semiconductor wafer to minimize dishing, while the high topography features are removed by polishing with a polishing pad and a polishing fluid.\nThe invention pertains to a process according to which, a resistant film is selectively formed on low topography features on a semiconductor wafer. Further, the invention pertains to a semiconductor wafer adapted for polishing with a polishing pad and a polishing fluid with which high topography features on the wafer chemically react and dissolve, and low topography features on the wafer have a resistant film that resists the polishing fluid. Further, the invention pertains to apparatus for forming a resistant film selectively on low topography features on a semiconductor wafer."}
{"text": "1. Field of the Invention\nThis invention relates to novel antibiotic compounds designated BU-3292T A and B. The compounds possess antimicrobial and antitumor activities.\n2. Description of the Prior Art\nElucidation of the structure of BU-3292T A revealed that it had a unique depsipeptide structure. This type of pentapeptide lactone has been proposed as a key intermediate in the biosynthesis of actinomycin (Katz. E: Antibiotics, Vol. II. Biosynthesis, p. 276, Springer-Verlay, Berlin, 1967), but none of such monomer type compound has so far been discovered in the microbial metabolites.\nActinomycin Z.sub.5 (Tetrahedron Letter, 2567, 1973) and Schering 1a and 1b (15th Interscience Conference on Antimicrobial Agents and Chemotherapy, Abstract No. 419) are reported to each contain one mole of 5-methylproline. BU-3292T B differs from these actinomycins in that it has two moles of 5-methylproline forming two identical peptide lactone side chains."}
{"text": "A novel gene MPR1 which was previously found in a yeast Saccharomyces cerevisiae S1278b strain by the present inventor (non-patent literature 1) codes for N-acetyltransferase (Mpr1) which detoxifies the toxic analogue of amino acid proline, azetidine-2-calboxylate (AZC), by the acetylation of AZC (FIG. 1) (non-patent literature 2).\nAZC enters into cells via proline permease and is incorporated into proteins by competing with proline upon protein synthesis. As a result of the incorporation of AZC, proteins with abnormal structures and impaired functions accumulate in cells and inhibit the cell growth. In cells which express MPR1, AZC is subjected to the N-acetylation in cytosol and is not incorporated into nascent proteins and therefore, such cells are thought to acquire AZC resistance (non-patent literature 3).\nAs a result of homology search, homologous genes which have similar functions as MPR1 from fission yeast Schizosaccharomyces pombe and, sibling species of S. cerevisiae, S. paradoxus (ppr1+, Spa MPR1) were identified (non-patent literatures 4 and 5) and therefore, MPR1 is believed to be widely distributed among yeasts. However, interestingly, S. cerevisiae S288C strain which was used for the genomic analysis of yeasts and sake yeast lack MPR1. MPR1 is located near the subtelomeric region of 14th chromosome of S1278b strain. A homologous gene (MPR2) which differs from MPR1 only at the 85th residue is located near the subtelomeric region of 10th chromosome and there is no difference in the functions between MPR1 and MPR2 (non-patent literature 1). Comparing S288C strain and S1278b strain, though some genes show base substitutions and deletions, no gene other than MPR1 that is present only in S1278b strain and codes for a certain phenotype has been reported.\nAZC is rare in nature and therefore, is probably not the intrinsic substrate of MPR1 product (Mpr1). Accordingly, the present inventors analyzed the biological function of Mpr1 and identified the intracellular substrate of the same (non-patent literature 6). Firstly, cells of MPR1•MPR2 disruptant were exposed to oxidative stresses such as hydrogen peroxide and heat shock. As a result, the MPR1•MPR2 disruptant showed decreased survival rate and increased intracellular reactive oxygen species (ROS) level compared to the wild-type strain. On the other hand, when multicopies of MPR1 were introduced into S. cerevisiae 288C strain which originally lacks MPR1 and MPR2, the strain exhibited increased survival rate under oxidative stress and decreased ROS level. Accordingly, it was confirmed that Mpr1 lowers the intracellular oxidative level. PUT2 disruptant which accumulates a metabolic intermediate of proline, Δ1-pyrroline-5-calboxylate (P5C), was confirmed to show growth impairment and increase in ROS level. Therefore, the role of Mpr1 in PUT2 disruptant was examined thereafter. The result showed that when MPR1 and MPR2 in PUT2 disruptant are disrupted, the growth is strongly inhibited and the ROS level is increased, whereas the overexpression of MPR1 decreased the ROS level. Further, the analysis using a recombinant enzyme showed that Mpr1 acetylates P5C or glutamate-gamma-semialdehyde (GSA) which is in equilibrium to P5C. From the above results, it is suggested that the intracellular accumulation of P5C leads to the generation of ROS which causes the cytotoxicity and that Mpr1 regulates the ROS level via the acetylation of P5C/GSA so as to alleviate the oxidative stress (FIG. 2) (non-patent literature 6).\nUnder the fermentative production environment, yeasts are subjected to various stresses such as cold temperature, freezing, drying, oxidation, high osmolarity, high ethanol concentration and biased nutrition. Long term exposure of yeast to such stresses as above brings about the cleavage of noncovalent bondings in the intracellular proteins which leads to exposure of hydrophobic amino acids on the surfaces of the proteins and generation of “abnormal proteins” which lost their normal structures and functions, and results in the restriction of the useful functions of the yeast. In particular, the oxidative stress is generated by various factors such as heat shock, hydrogen peroxide, freezing (freezing-thawing) damage, high ethanol concentration and the like and is a great cause of the inhibition of the growth of the yeast cells. In the fields of fermented food products and brewed food products, there is a desire for breeding a yeast which is highly resistant to the oxidative stress.\nSo far, the present inventors found that an amino acid proline has a property of protecting yeasts from stresses such as freezing, drying, oxidation and the like (Patent literature 1). Further, the present inventors found that a yeast strain whose gene encoding proline degradative enzyme was disrupted by the genetic engineering acquires the ethanol resistance by accumulating proline in the cells (Patent literature 2).    Patent literature 1: JP-A-9-234058    Patent literature 2: JP-A-2006-67806    Non-patent literature 1: H. Takagi et al., J. Bacteriol., 182, 4249-4256 (2000)    Non-patent literature 2: M. Shichiri et al., J. Biol. Chem., 276, 41998-42002 (2001)    Non-patent literature 3: C. Hoshikawa et al., Proc. Natl. Acad. Sci. U.S.A., 100, 11505-11510 (2003)    Non-patent literature 4: Y. Kimura et al., Yeast, 19, 1437-1445 (2002)    Non-patent literature 5: M. Nomura et al., J. Biochem., 133, 67-74 (2003)    Non-patent literature 6: M. Nomura et al., Proc. Natl. Acad. Sci. U.S.A., 101, 12616-12621 (2004)"}
{"text": "As an electric current type of light-emitting device, organic light emitting diodes (OLEDs) have found wide application in high performance displays. With the increase in the size of the display, traditional passive matrix organic light-emitting diode (PMOLED) displays require shorter driving time for a single pixel, which necessitates an increased transient current and increased power consumption. Meanwhile, application of a large current results in a large voltage drop on the ITO wire and a high working voltage of the OLED, and in turn decreases its efficiency. Active matrix organic light-emitting diode (AMOLED) displays may address these issues elegantly by means of switch transistors scanning and inputting currents for OLEDs line by line.\nAmong the three types of driving approaches for AMOLED, which are digital driving, current driving and voltage driving, the voltage driving approach is similar to the traditional driving approach for active matrix liquid crystal displays (AMLCDs), i.e., providing by a driving chip (IC) a voltage signal that represents a grayscale, which voltage signal would be converted inside a sub-pixel into a current signal for driving a thin film transistor, so as to drive the OLED to exhibit a luminance grayscale. This approach is advantageous in that it is fast in driving speed, simple to implement, and appropriate for driving large-sized panels, and thus is extensively applied in the industry.\nHowever, in an AMOLED display of the voltage driving type, the current-luminance (I-L) conversion efficiency decreases as the OLED ages over time. Even if the currents are the same, the luminance displayed may be different because due to their different aging degrees the OLEDs are different in their conversion efficiencies. This leads to an issue that non-uniformity of the luminance is present in the images displayed on the display panel.\nIn order to improve the luminance uniformity, external compensation approaches are generally applied at present for the AMOLED display of the voltage driving type. Namely, each sub-pixel of the display panel is connected with a driving chip through a respective sense line corresponding one-to-one thereto, which driving chip detects the aging of the OLEDs in respective sub-pixels through the respective sense lines and then performs compensation of the sub-pixels in accordance with the detection. Nevertheless, in the above display panel, each sub-pixel is connected to a respective sense line, leading to increased wirings in the display panel, which is disadvantageous for fabrication of a high resolution display panel. Moreover, the number of the signal channels of the driving chip will be doubled, resulting in an increased area of the driving chip and a high cost."}
{"text": "The present invention relates to an orientation-adjusting device, and more particular to an orientation-adjusting device for properly adjusting the orientation of a wireless transceiver so as to allow the transceiver to be fitted into a confined space.\nUniversal Serial Bus (USB) has been becoming a major type of connecting ports for personal computers. As result, many newly developed peripheral equipment, e.g. Bluetooth-related peripheral equipment, are designed with USB as a standard interface. On the other hand, a wireless transceiver is more and more popular to become a part of a personal computer. Therefore, the equipment of a personal computer with a wireless transceiver interfaced therebetween with a USB adapter becomes a trend. In order to be fitted into a confined space, especially in a portable computer of small size, and perform the wireless communication function well, the wireless transceiver should be properly oriented.\nTherefore, an object of the present invention is to provide an orientation-adjusting device for properly orienting a wireless transceiver in a confined space, which preferably performs two dimensional adjustment.\nA first aspect of the present invention relates to an orientation-adjusting device for adjusting an orientation of a wireless communication device. The wireless communication device is mounted in a main frame via an interface device, and the orientation-adjusting device includes a first housing for accommodating therein the wireless communication device; a second housing pivotally connected to the first housing to allow the first housing to rotate in a first direction relative thereto; and a third housing for accommodating therein the interface device, pivotally connected to the second housing to allow the second housing to rotate in a second direction relative thereto.\nPreferably, the second housing includes a passage for penetrating therethrough a signal cable connecting the wireless communication device and the interface device.\nIn an embodiment, the first housing includes a first cylindrical portion for sleeving therearound the second housing, and serving as a first shaft for rotating the first housing therewith in the first direction.\nIn an embodiment, the second housing includes a first engaging part, the third housing includes a second engaging part, and the first and second engaging parts pivotally connected to each other to serve as a second shaft for rotating the second housing therewith in the second direction, and form a second hollow cylindrical portion for penetrating therethrough the signal cable.\nIn an embodiment, the second housing includes a hollow disk body secured to the first engaging part for sleeving around the first cylindrical portion. Preferably, the second hollow cylindrical portion serves as a rotation-angle confiner of the first housing in the first direction.\nPreferably, the first direction is perpendicular to the second direction so as to perform two-dimensional adjustment.\nFor example, the wireless communication device can be a wireless transceiver, the main frame can be a personal computer, and the interface device can be a Universal Serial Bus (USB) adapter.\nA second aspect of the present invention relates to another orientation-adjusting device for adjusting an orientation of a wireless communication device. The orientation-adjusting device includes a first housing for accommodating therein the wireless communication device, which includes a first shaft; a second housing including a first portion penetrating therethrough the signal cable and sleeving around the first shaft, and a second portion being of a hollow cylindrical shape; and a third housing including a third portion for accommodating therein the interface device, and a fourth portion being of a hollow cylindrical shape, engaging with the second portion of the second housing to form a hollow cylinder for penetrating therethrough the signal cable, and serving as a second shaft, wherein the first and second housings optionally perform a first relative rotating motion with the first shaft, and the second and third housings optionally perform a second relative rotating motion with the second shaft to adjust the orientation of the wireless communication device.\nIn an embodiment, the first portion of the second housing being of a disk shape.\nPreferably, the first and second relative rotating motions are in different directions so as to perform two-dimensional adjustment.\nA third aspect of the present invention relates to an orientation-adjustable transceiver assembly, which includes a transceiver mounted in a first housing; a signal cable connected to the transceiver and penetrating through a second housing which pivots relative to the first housing to adjust the orientation of the transceiver in a first direction; and an interface device connected to the signal cable and mounted in a third housing which pivots relative to the second housing to adjust the orientation of the transceiver in a second direction different from the first direction.\nIn an embodiment, the first housing includes a first shaft for sleeving therearound the second housing, and allowing the second housing to pivot thereabout in the first direction.\nIn an embodiment, the second housing includes a first engaging part, the third housing includes a second engaging part, and the first and second engaging parts engage with each other to form a second shaft for allowing the third housing to pivot thereabout in the second direction.\nPreferably, the second shaft is hollow for penetrating therethrough the signal cable.\nPreferably, the interface device is a Universal Serial Bus (USB) adapter."}
{"text": "This invention relates generally to human prostheses, and especially to spinal column vertebral disc prostheses. The invention also relates to surgical procedures for preparing the patient to receive a vertebral disc endoprosthesis, and for implanting that endoprosthesis in the patient's spine.\nThe herniation of a spinal disc and the often resultant symptoms of intractable pain, weakness, sensory loss, incontinence and progressive arthritis are among the most common of debilitating processes affecting mankind. If a patient's condition does not improve after conservative treatment, and if clear physical evidence of nerve root or spinal cord compression is apparent, and if correlating radiographic studies (i.e., MRI or CT imaging or myelography) confirm the condition, surgical removal of the herniated disc may be indicated. The process of discectomy--as the name implies--involves the simple removal of the disc without attempt to replace or repair the malfunctioning unit. In the United States in 1985, over 250,000 such operations were performed in the lumbar spine and in the cervical spine.\nStatistics suggest that present surgical techniques are likely to result in short-term relief, but will not prevent the progressive deterioration of the patient's condition in the long run. Through better pre-operative procedures and diagnostic studies, long-term patient results have improved somewhat. But it has become clear that unless the removed disc is replaced or the spine is otherwise properly supported, further degeneration of the patient's condition will almost certainly occur.\nIn the mid-1950's and 60's, Cloward and Smith & Robinson popularized anterior surgical approaches to the cervical spine for the treatment of cervical degenerative disc disease and related disorders of the vertebrae, spinal cord and nerve root; these surgeries involved disc removal followed by interbody fusion with a bone graft. It was noted by Robinson (Robinson, R. A.: The Results of Anterior Interbody Fusion of the Cervical Spine, J. Bone Joint Surg., 440A: 1569-1586, 1962) that after surgical fusion, osteophyte (bone spur) reabsorption at the fused segment might take place. However, it has become increasingly apparent that unfused vertebral segments at the levels above and below the fused segment degenerate at accelerated rates as a direct result of this fusion. This has led some surgeons to perform discectomy alone, without fusion, by a posterior approach in the neck of some patients. However, as has occurred in surgeries involving the lower back where discectomy without fusion is more common as the initial treatment for disc herniation syndromes, progressive degeneration at the level of disc excision is the rule rather than the exception. Premature degenerative disc disease at the level above and below the excised disc can and does occur.\nSpine surgery occasionally involves fusion of the spine segments. In addition to the problems created by disc herniation, traumatic, malignant, infectious and degenerative syndromes of the spine can be treated by fusion. Other procedures can include bone grafts and heavy duty metallic rods, hooks, plates and screws being appended to the patient's anatomy; often they are rigidly and internally fixed. None provide for a patient's return to near-normal functioning. Though these procedures may solve a short-term problem, they can cause other, longer term, problems.\nA number of attempts have been made to solve some of the problems described above by providing a patient with spinal disc prostheses, or artificial discs of one sort or another. For example, Steffee, U.S. Pat. No. 5,031,437, describes a spinal disc prosthesis having upper and lower rigid flat plates and a flat elastomeric core sandwiched between the plates. Frey et al., U.S. Pat. Nos. 4,917,704 and 4,955,908, disclose intervertebral prostheses, but the prostheses are described as solid bodies.\nU.S. Pat. Nos. 4,911,718 and 5,171,281 disclose resilient disc spacers, but no inter-connective or containing planes or like elements are suggested, and sealing the entire unit is not taught.\nIt is the primary aim of the present invention to provide a vertebral disc endoprosthesis which will perform effectively and efficiently within a patient's spine over a long period of time, and which will not encourage degeneration of or cause damage to adjacent natural disc parts.\nIt is a related objective to provide a new vertebral disc endoprosthesis surgical procedure which will decrease post-operative recovery time and inhibit post-operative disc, vertebral body and spinal joint degeneration.\nIt is yet another object to provide a method of installing the endoprosthesis so as to accurately mate the endoprosthesis with an adjacent specifically formed bone surface. An associated object is to provide an endoprosthesis which will encourage bone attachment to, and growth upon, adjacent outer surfaces of the endoprosthesis.\nYet another object is to provide a vertebral endoprosthesis in which the parts are non-oncogenic.\nStill another object is to provide a vertebral disc endoprosthesis having a resilient element to accommodate shocks and other forces applied to the spine.\nAnother object is to provide a highly effective vertebral endoprosthesis which includes several disc endoprostheses and one or more prosthetic vertebral bodies. A related object is to provide these elements in a pre-assembled array for implantation in a patient."}
{"text": "Endodontics or root canal therapy is that branch of dentistry that deals with the diseases of the dental pulp and associated tissues. One aspect of endodontics comprises the treatment of infected root canals, the removal of diseased pulp tissues, followed by the biomechanical modification and the subsequent filling of the pulp canal (root canal). Root canal therapy is generally indicated for teeth having sound external structures but having diseased, dead or dying pulp tissues. Such teeth may or may not generally possess intact enamel and dentin and are satisfactorily engaged with bony tissue. In such teeth, the pulp tissue and excised portions of the root should be replaced by a biocompatible substitute. One technique for the preparation of a root canal involves creating a coronal access opening with a conventional dental drill. A tool is used for gross removal of pulp material from the root canal through the coronal access opening. The void formed is enlarged with reamers and/or files to result in a fully excavated cavity. Debris is removed from this cavity by flushing and the cavity is cleansed to remove all diseased tissue. Following chemical antisepsis, the excavated canal is ready for filling.\nA basic method involves inserting a filling cone into a root canal and cementing therein to obturate the canal. The common root canal filling cone material is made from gutta percha. Lateral condensation is a method in which several filling cones, a primary cone and auxiliary cones, are inserted into a root canal. The primary cone is inserted and cemented to the seat of the root canal. Using a tapered spreader, the primary cone is then squeezed against the side of the root canal and a second cone is inserted and cemented into place. This process is continued until the root canal is completely obturated which can require up to 10 to 15 filling cones. Vertical condensation of warm or hot gutta percha is yet another method of sealing root canals. After cementing a primary cone short of the apex of the root canal, heat application is alternated with a series of smaller and smaller pluggers until the gutta percha is moved to the apex. This is often possible when the smallest plugger approaches the apex of the tooth within 3 to 5 millimeters. The space is then backfilled. Lateral canals are packed and sealed as a consequence of lateral expansion of a wave of heated gutta percha. Alternatively, small segments of gutta percha can be used in this method that are inserted into the root canal, heated in order that they can adhere to one another and each backfilled one at a time until the root canal is filled. All three of these methods, the single filling cone, lateral condensation and vertical condensation apply root canal cement or sealer around the individual cones or in between segments as a binding agent.\nAnother method employs an injection gun that injects warm or hot gutta percha filling material into a root canal. The injector initially places heated gutta percha at the apical area of the root canal through a needle-like canula tip and fills the gutta percha into any surrounding voids/spaces under pressure or at the seat of the root canal which is then condensed with a plugger into the root tip. The injector then backfills the root canal by injecting additional gutta percha into the root canal until it is obturated. A similar method involves heating gutta percha on a flexible metal or plastic carrier used to insert the gutta percha into the root canal. The carrier may be a solid rod, or a hollow rod, situated in the center of a master cone. The rod is connected to a handle which may be removed by slipping it out of the hollow rod, or cutting it off if it is a solid rod.\nMost of the current methods employed in obturating a canal use a gutta percha material that is inert in nature and will not be absorbed or degraded by the living tissue if the root canal is overfilled and extends beyond the apex. It has been a challenge for dentists to control the exact amount of the material within the border of the root canal to avoid overfilling. Moreover, gutta percha material is a polyisoprene rubber material in nature, which does not have the capability to bond to most of the dental materials, especially when the root canal sealer is a polymer-based material. Gutta percha exhibits poor strength and brittleness. Dental gutta percha points/cones tend to break in harsh conditions, e.g., sharply curved root canals, tight spaces during a root canal treatment, and the like.\nIt is desirable to provide a root canal filling material that bonds easily to sealants. It is preferable that the root canal filling material have proper strength and flexibility. It would be beneficial that the root canal filling material be bioactive."}
{"text": "“Web 1.0” is the term associated with the first generation of internet browser applications and programs, along with the associated client-side software entities and server-side software entities used to support and access information using the Internet. Such Web 1.0 technologies, like most first-generation technologies, are geared more to enabling a workable system and to the capabilities of the available software and hardware platforms, rather than to creating a rich and efficient experience for the system's users. Thus, conventional Web 1.0 technologies, while efficient for machines, are often highly inefficient and frustrating for their human users.\nIn particular, Web 1.0 technologies operate on a “click-wait” or a “start-stop” philosophy. That is, when a user wishes to view a web page, the user must generate a request using the client-side browser software, and send that request to the server. The user must then wait for the server to respond to the request and forward the requested data. The user must further wait for all of the requested data to be received by the client-side browser software and for the browser software to parse and display all of the requested information before the user is allowed to interact with the requested web page.\nThis is frustrating for most users on a number of levels. First, for slow or bandwidth-limited Internet connections, obtaining all of the requested data can often take a relatively long time. Furthermore, even when the user has high-speed access to the Internet, a web page that requires data to be re-loaded or refreshed on a fairly regular basis, such as mapping web pages, sporting events scores, or play-by-play web pages and the like, can cause significant delays. This is typically due to Web 1.0 requirements that the entire web page be retransmitted even if no or only minimal changes have occurred to the displayed information.\nAccordingly, the next generation of technologies used to access and support the Internet is currently being developed and collected under the rubric “Web 2.0”. A key feature in the “Web 2.0” concept is to eliminate the above-outlined “click-wait” or “start-stop” cycle, by asynchronously supplying data associated with a particular web page to the user from the associated web server. The transfer occurs as a background process, while a user is still viewing and possibly interacting with the web page, which anticipates the fact that the user will wish to access that asynchronously-supplied data. A number of important technologies within the “Web 2.0” concept have already been developed. These include “AJAX”, SVG, and the like.\nAsynchronous JavaScript and XML, or “AJAX”, is a web development technique used to create interactive web applications. AJAX is used to make web pages feel more responsive by exchanging small amounts of data between the client application and the server as a background process. Accordingly, by using AJAX, an entire web page does not have to be re-loaded each time a portion of the page needs to be refreshed or the user makes a change to the web page at the client side. AJAX is used to increase the web page's interactivity, speed, and usability. AJAX itself makes use of a number of available techniques and technologies, including XHTML (extended hypertext markup language) and CSS (cascading style sheets), which are used to define web pages and provide markup and styling information for the web pages. It also makes use of a client-side scripting language, such as JavaScript, that allows the DOM (document object model) to be accessed and manipulated, so that the information in the web page can be dynamically displayed and can be interacted with by the user.\nOther important technologies include the XMLHttpRequest object, which is used to exchange data asynchronously between the client-side browser software and the server supporting the web page being displayed, and XML, RSS and other data exchange standards, which are used as the format for transferring data from the server to the client-side browser application. Finally, SVG (scalable vector graphics) is used to define the graphical elements of the web page to be displayed using the client-side browser application.\nIn addition to Web 1.0 and Web 2.0 technologies, an entirely different set of software technologies are used to access other data available over local area networks, wide area networks, the internet and the like. These technologies are traditionally referred to as “client-server applications”, where a complex software application having a rich set of features is installed on a particular client computer. This software application executes on the client computer and is used to access, display and interact with information stored on a server that is accessed via a local area network, a wide area network, the Internet or the like. While such client-server applications allow for dynamic displays and make manipulating information easy, such client-server applications are difficult to deploy to all of the client machines, and are difficult to update.\nTimelines for displaying data are known. The SIMILE project has a thin client Timeline that is interactive. However, it lacks many desirable features. For example, there is no toolbar that is able to interact with the Timeline. Also, the SIMILE timeline only supports a proprietary data format and does not ingest standard formats like RSS and GeoRSS. Once the data is read in, it does not link to other visualizations or support selection or highlighting. The ability to zoom in or out and the ability to change scales are also not included in SIMILE's Timeline.\nThe focus of “Business Intelligence” is to enable corporations to analyze their business data and gain insights that result in better business performance. Business Intelligence software vendors have been successful by providing better reporting tools, analytic dashboards, and data analysis capability that help corporations improve operating efficiency. One type of data that is underutilized by Business Intelligences systems is that involving locations. It is estimated that 80% of all data maintained by corporations has a location component. Location data is either stored directly, such as storing the latitude and longitude of tracked objects in a database, or the location is implied, such as storing the street address which can be decoded into geo-coordinates. In the past, geographic data was passed to the GIS (Geographic Information System) department for processing where reports were prepared for the analysts and decision makers. These GIS systems were heavy desktop applications, which needed dedicated machines and highly trained staff to use. This process was often lengthy, and therefore rarely used by analysts needing to make quick decisions.\nLocation data is traditionally latitude-longitude based data. Map companies, for example, present information in the form of latitude-longitude based data. Items that have map based coordinates may include merchandise shelf position or aisle location in a store.\nPlanograms are also a tool used by retailers to optimize the layout and location of merchandise. In this application the planogram is a visual planning tool used by merchants and planners to analyze, plan and communicate to stores how product should be laid out and configured. Retailers utilize this application of the planogram more than consumer goods manufacturers. For the most part, retailers guard their planograms closely as competitive information.\nRecently, with advances in hardware technology and the emergence of “Web 2.0 AJAX” applications, these untapped location data sources can now be made available to users. Thin client, geospatial mapping platforms provide innovative GIS functionality. These platforms are inexpensive, are bundled with image and map content, provide rich functionality, and are totally browser-based. Today, it is possible to do live interactive maps without needing to install complex client software. As a result they have dramatically broadened access to GIS and mapping technology.\nIn parallel with the emergence of these thin client GIS web platforms, new standards have emerged for publishing and consuming geo-encoded information. Three standards are widely used. First, GeoRSS is a rather simple extension of RSS, an XML specification for publishing information best described as a sequence of items in channels. GeoRSS extends RSS by including geographical encodings to describe the location of the event items. These encodings can range from a simple Latitude/Longitude point location up to a complex polygon describing a geographical region. These points are encoded using the GML (geographic markup language) standard. Second, KML (Keyhole Markup Language, named after the company which developed it, Keyhole Technologies) is the language used to describe three dimensional spatial planes. Third, GPX is a standard popularized by Garmin GPS units to encode “GPS Tracks” of the device. These formats are similar enough that many open source programs exist to translate between the three data formats.\nWeb applications that combine geospatial maps with geo-encoded information from various sources are often called mashups. Mashups typically take spatial and non-spatial data, combining or “mashing” the data sets together to a display. Common examples of mashups are “Amazon book purchases overlaid by region,” “Real estate price by neighborhood,” and “earthquake events published by the USGS” plotted geospatially on a Google Map. Although the public domain data was always available, the merging of the data onto a map was difficult. What has changed is that good-enough GIS software is widely available and interesting geo-encoded datasets are now published using geospatial standards. As a result, the number of mashup applications has grown dramatically and there are literally hundreds of thousands of mashup sites."}
{"text": "As large native substrates for group III-Nitride (III-N) semiconductors are not yet widely available, III-N films, such as GaN and its alloys, are currently grown by heteroepitaxy on suitable non-III-N substrates. Typically, the films are grown on sapphire (Al2O3), silicon carbide (SiC), or silicon substrates. Silicon substrates are emerging as a particularly attractive substrate candidate for III-N layers due to their low cost, wide availability, large wafer sizes, thermal properties, and ease of integration with silicon-based electronics. However, due to the large lattice mismatch and thermal expansion coefficient mismatch between silicon and III-N materials, there is typically a net tensile stress in III-N epitaxial layers deposited directly on silicon substrates.\nThis mismatch can result in cracking of the layers and thus thick III-N layers that are crack-free and that exhibit adequate structural quality can be difficult to grow on silicon substrates. For example, the maximum thickness of III-N layers that can be grown without sustaining substantial defects may be limited. If the III-N layers are grown too thick, tensile stress that develops in the layers upon cooling becomes substantial, which can cause cracking of the layers. In many applications in which III-N heteroepitaxial layers are used, it may be necessary that very thick III-N layers of adequate quality be grown on the foreign substrates."}
{"text": "The present invention relates to ink pads for use in self-inking rubber stamps and pre-inked stamps both of which are well known in the art. The prior art originally taught an ink stamp pad separate from a rubber stamp to ink the stamp. These conventional ink pads varied in size and were typically contained in plastic or metal containers with a closure lid to prevent the ink from drying out when the pad was not in use. Those pads, generally referred to as conventional ink pads, were usually formed of cotton felt and then enclosed with a thin fabric of cotton or linen.\nThe technology of rubber stamps moved on to self-inkers, such as date stamps, where the ink pad was contained within the stamp, and the rubber stamp was held against the ink pad until the stamp was used whereupon the rubber stamp through a cam follower would rotate 180 degrees to a downward direction into engagement with the medium being stamped. The ink pads of these self-inkers are contained in an open topped tray and the pad is composed of a foam material having an open cell micropore structure which contains and holds the ink. Attached to the top surface of the foam pad is a thin layer of suede-like fabric which provides the necessary uniform ink transfer from the foam pad to the rubber stamp. Foam ink pads can only be used with water-based inks since oil and alcohol based inks will cause the pad to swell and render it unusable. For that reason, foam-type ink pads cannot be used with fast drying inks which are alcohol based.\nWhen the self-inkers are not in use, the rubber stamp is held against the pad with a spring force which deforms the soft foam pad to the point the inkpad takes an impression of the stamp. This fixed impression also causes the edges of the ink pad to curl upward above the top edge of the tray and create problems when the tray is removed or inserted into the stamp.\nPre-inked stamps which are another modern concept include a layer of material (EVA) which is flashed with light with the image of the rubber stamp blocking the light so that the image of the stamp on the EVA material is open to pass ink onto the medium being stamped. The portion of the EVA material which light contacts causes the material to seal so the ink cannot pass through while the portion blocked from light by the stamp image freely passes ink to the medium being stamped. This pre-inked stamp process eliminates the step of inking the rubber stamp before the image is stamped on the medium. The ink cartridges of these pre-inked stamps utilize a proprietary foam material.\nThe ink pad of the present invention is formed from a non-woven polypropylene filament which is compressed and heated to a substantially rigid form. The ink pad can be used with all types of inks since alcohol and oil based inks do not effect the polypropylene filament. Attached to the ink pad is a second layer of suede-like cloth which is bonded by a thermoplastic adhesive to form a unitary structure. The two-layer ink pad is contained within an open top tray as are the foam pads in the prior art and the trays fixedly position the pad within the overall stamp frame structure.\nThe tightly woven suede-like material has a relatively short nap and provides a good even ink transfer to the rubber stamp. Since the overall ink pad is relatively rigid, there are no problems with the edges of the pad turning up as with the foam pads. The compressed non-woven polypropylene filament absorbs as much ink as the form pads of the prior art and releases the ink in an acceptable manner when the stamp is put in use. The pads are made in large sheets that are readily cut to size by knife edges in an automated process well known in the prior art. The pad of the present invention can also be used as a conventional ink stamp pad for rubber stamps as mentioned above.\nIt is therefore the principal object of the present invention to provide an ink pad for conventional stamp pads, self-inkers, and pre-inked stamps which is relatively inexpensive to construct.\nAnother object of the present invention is to provide an ink pad which can be used with water, oil and alcohol based inks.\nA further object of the present invention is to provide an ink pad which is substantially rigid while containing ample amounts of ink with superior ink release characteristics."}
{"text": "1. Field of the Invention\nThe present invention relates to an electron beam lithography apparatus using a patterned emitter. More particularly, the present invention is directed to an electron beam lithography apparatus using a pyroelectric plate on which a patterned metal thin layer is formed as an electron beam source.\n2. Description of the Related Art\nConventional apparatus for performing ferroelectric switching lithography using a patterned emitter emit electrons by switching a patterned ferroelectric emitter to expose an electron resist on a substrate to form a desired pattern which is the same as that of the emitter. The conventional ferroelectric switching emission, however, has a drawback in that an electrode formed on the emitter by a mask absorbs electrons. Moreover, the emitter cannot reliably emit electrons when it is not connected to the electrode.\nA lithography apparatus, which has a thermal electron source, controls an electron beam using an electro-magnetic field. In order to deploy a large-sized electron source, however, the lithography apparatus is required to include a device for collimating radially diffused electron beams. This requirement introduces complications in the design and structure of the lithography apparatus.\nA photocathode projection type lithography apparatus deploys a large-sized electron source but is extremely sensitive to air pollution and therefore cannot be commercially used.\nThe present invention provides an electron beam lithography apparatus capable of emitting electrons using an electron beam source without applying a high voltage.\nThe present invention also provides a method of fabricating an electron lithography apparatus that heats an electron beam source at a low temperature to emit electrons, and a method of making the electron beam source.\nAccording to an embodiment of the present invention, there is provided an electron beam lithography apparatus for providing a one-to-one projection of a pattern, comprising a pyroelectric emitter which is disposed from a substrate holder at a predetermined distance, the emitter including a pyroelectric plate having a patterned dielectric layer on its surface facing the substrate holder, wherein portions of the surface of the pyroelectric plate exposed by the patterned dielectric layer is covered with a metal thin layer, a heating source which heats the emitter, and a permanent magnet or an electro-magnet disposed outside the emitter and the substrate holder so as to control the paths of electrons emitted by the emitter.\nPreferably, the metal thin layer is formed on the pyroelectric plate that is exposed by the patterned dielectric layer. An adhesion layer of a predetermined thickness may be formed between the pyroelectric plate, and the dielectric layer and the metal thin layer. The dielectric layer and the pyroelectric plate may be sequentially formed below the metal thin layer.\nAccording to another embodiment of the present invention, there is provided an electron lithography apparatus for providing a x-to-one projection of a pattern, comprising a pyroelectric emitter which is disposed from a substrate holder at a predetermined distance, the emitter including a pyroelectric plate having a patterned dielectric layer on its surface facing the substrate holder, wherein portions of the surface of the pyroelectric plate exposed by the patterned dielectric layer is covered with a metal thin layer, a heating source which heats the emitter, and a deflector unit which is disposed between the emitter and the substrate holder so as to control the paths of electrons emitted by the emitter.\nAccording to yet another embodiment of the present invention, there is provided an electron beam lithography apparatus for providing a one-to-one projection of a pattern, comprising a pyroelectric emitter which is disposed from a substrate holder at a predetermined distance, the emitter including a pyroelectric plate having a patterned metal thin layer on its surface facing the substrate holder, wherein a dielectric layer is formed on the pyroelectric plate below the patterned metal thin layer, a heating source which heats the emitter, and a permanent magnet or an electro-magnet disposed outside the emitter and the substrate holder so as to control the paths of electrons emitted from the emitter.\nAccording to still another embodiment of the present invention, there is provided an electron beam lithography apparatus for providing a x-to-one projection of a pattern, comprising a pyroelectric emitter which is disposed from a substrate holder at a predetermined distance, the emitter including a pyroelectric plate having a patterned metal thin layer on its surface facing the substrate holder, wherein a dielectric layer is formed on the pyroelectric plate below the patterned metal thin layer, a heating source which heats the emitter, and a deflector unit which is disposed between the emitter and the substrate holder so as to change the paths of electrons emitted from the emitter.\nAccording to a feature of the present invention, the heating source may be a contact-type heating plate using a resistive-type heating or a remotely controlled heater that generates infrared rays.\nAccording to another feature of the present invention, the pyroelectric plate may be formed of a pyroelectric material selected from the group consisting of LiNbO3, LiTaO3, and BaTiO3.\nAccording to yet another feature of the present invention, the dielectric layer or the adhesion layer may be formed of a dielectric material selected from the group consisting of SiO2, Si3N4, and TiO2.\nAccording to still another feature of the present invention, the metal thin layer may be formed of metal selected from the group consisting of Ti, Au, Pt, Ta, and Al.\nAccording to another feature of the present invention, the deflector unit may include deflection plates which deflect electrons emitted from the emitter, at least one magnetic lens which collimates the deflected electrons, and an aperture which filters the collimated electrons.\nAccording to yet another feature of the present invention, the heating source may heat the emitter to a temperature up to 120xc2x0 C."}
{"text": "1. Field of the Invention\nThe invention relates to a stirrup.\n2. Description of the Prior Art\nA stirrup consists of a pair of curved arms which serve as a carrier for an eye and a foot. Each arm is joined to the other at one end while the opposite ends of the arms are spaced from one another. The eye is located at the joint and the foot is disposed between the spaced ends of the arms.\nThe stirrup is secured to a horse by passing a strap through the eye which then serves as a holding member for the strap. The arms and the foot can rotate relative to the eye between a position in which the eye is coplanar with the arms and positions in which the eye is transverse to the arms. This allows the eye to be oriented such that the strap is flat against the horse and flat against the leg of a rider. At the same time, the arms and the foot can be properly positioned for the foot of the rider.\nThe foot of the stirrup can rotate relative to the arms on an axis extending between the spaced ends of the arms. The foot is rotatable forwards and backwards from a central position in which a plane defined by the arms is perpendicular to the tread of the foot. The primary reason for rotation of the foot is safety. Thus, should the rider fall off the horse backwards, the foot of the stirrup can rotate backwards to allow the foot of the rider to be released from the stirrup more easily.\nThe stirrup described above has a drawback related to the eye. In this stirrup, the eye has several discrete positions relative to the arms and the foot. When the eye is moved into one of these positions, the eye is locked to the arms and the foot by a locking mechanism. Accordingly, the eye rotates when the foot of the rider turns or twists sideways in the stirrup. As the eye rotates, the strap passing through the eye twists so that the edges of the strap are turned towards and rub against the horse and the leg of the rider. This is uncomfortable for both the horse and the rider.\nAnother drawback of the preceding stirrup stems from the fact that the foot of the stirrup can rotate forwards as well as backwards from the central position. The foot of the rider tends to push forwards with a resultant forward rotation of the foot of the stirrup. This causes the rider to experience a certain degree of instability."}
{"text": "Field of the Invention\nThe present invention relates to a technique for calculating feature amounts from an image.\nDescription of the Related Art\nA method for searching for a similar image using local feature amounts (local feature amount) of an image has been proposed. In this method, first, a characteristic point (a local feature point) is extracted from an image (for example, refer to C. Harris and M. J. Stephens, “A combined corner and edge detector,” In Alvey Vision Conference, pages 147-152, 1988). Next, based on that local feature point and image information of a periphery thereof, a feature amount (a local feature amount) corresponding to that local feature point is calculated (for example, refer to David G. Lowe, “Distinctive Image Features from Scale-Invariant Keypoints,” International Journal of Computer Vision, 60, 2 (2004), pp. 91-110).\nIn a method using local feature amounts, a local feature amount is defined as information comprised by a plurality of elements that are rotationally invariant and magnification/reduction invariant in a plane. With this, a search is made to be possible even in a case where an image is rotated, magnified or reduced in a plane. Generally, a local feature amount is represented as a vector.\nIn David G. Lowe, “Distinctive Image Features from Scale-Invariant Keypoints,” International Journal of Computer Vision, 60, 2 (2004), pp. 91-110, for example, to extract a local feature amount that is rotationally invariant in a plane, a main direction is calculated from a pixel pattern of a local region of a local feature point periphery, and a direction normalization is performed by rotating the local region in the plane based on the main direction when calculating the local feature amount. Also, in order to calculate a magnification/reduction invariant local feature amount, an image of a different scale is generated internally, and from each image scale, local feature point extraction and local feature amount calculation are performed. Here, an image collection of a sequence of differing scales generated internally is generally called a scale space.\nBy the foregoing method, a plurality of local feature points are extracted from one image. In an image search using local feature amounts, matching is performed by performing a comparison of local feature amounts calculated from respective local feature points with each other. In a voting method (for example, refer to Japanese Patent Laid-Open No. 2000-57377), which is often used, if there exists in a registered image (hereinafter referred to as a comparison target image) a feature point having a feature amount similar to the local feature amount of a feature point that is extracted from a comparison source image, a vote is made for the comparison target registered image. Configuration is such that the greater the number of votes for a registered image, the more similar it is to the comparison source image.\nThere was a problem in that upon performance of a search of an image subjected to an out-of-plane rotation using local feature amounts as recited in David G. Lowe, “Distinctive Image Features from Scale-Invariant Keypoints,” International Journal of Computer Vision, 60, 2 (2004), pp. 91-110, while a search of an image subjected to an out-of-plane rotation of up to 30 degrees is possible, the search accuracy for an image subjected to an out-of-plane rotation greater than that would become low.\nFIGS. 1A-1D are views for explaining an example of a pixel range used for calculating a local feature amount. An explanation is given for the cause of lower search accuracy of an image which is subjected to an out-of-plane rotation using FIGS. 1A-1D. FIG. 1A is a schematic diagram of the image captured from a front surface of an object, and FIG. 1B is a schematic diagram of the image captured from the same location after causing the object to rotate 60 degrees counterclockwise from the state of FIG. 1A. Here, reference numeral 101, whose center point is one of local feature points extracted from the image of FIG. 1A, illustrates an example of a pixel range used for calculating a local feature amount from the local feature point. Similarly, reference numeral 102, whose center point is one of local feature points extracted from the image of FIG. 1B, is illustrating an example of a pixel range used for calculating a local feature amount from the local feature point. Also, FIG. 1C is a figure magnifying a periphery of the pixel range 101 in FIG. 1A and FIG. 1D is a figure magnifying a periphery of the pixel range 102 in FIG. 1B.\nHere, 103 and 104 are local feature amount calculation regions whose center is a local feature point extracted from a bottom-left corner of a character “B”. The left side of the local feature amount calculation region 103 is slightly overlapping the character “A” and the right side is substantially in the middle of the character “B”. On the other hand, the left side of the local feature amount calculation region 104 is overlapping more than the half of the character “A” and the right side is the same as the right side of “B”. In spite of the fact that the local feature points are extracted from almost the same location in the local feature amount calculation region 103 and the local feature amount calculation region 104, it can be seen that the pixels used for calculating the local feature amounts are significantly different. With this, the local feature amount changes and the precision of an image search is degraded. In this way, it was conventionally difficult to search an image subjected to an out-of-plane rotation with a high precision with local feature amounts that are rotationally invariant in a plane."}
{"text": "1. Field of the Invention\nThe present invention relates generally to redundancy repair of integrated circuit memory devices, and more particularly to a memory device having row redundancy repair.\n2. Brief Description of the Related Art\nElectronic memory, including read only memory (ROM), flash memory, and random access memory (RAM) devices, such as DRAMs and SRAMs, typically utilize memory arrays having rows and columns of memory cells for storing data. Manufactured memory arrays are tested for integrity, and defects are repaired. Typical tests involve writing data to the array using an external pattern, and reading the pattern back to confirm that the data is stored accurately.\nWhen testing reveals a defective primary row or column, a non-defective redundant row or column is used to replace the defective primary row or column. This is achieved with fuses or anti-fuses that set hardwired latches associated with the defective row or column and map the address of the defective primary row or column to a fully-operable redundant row or column. With this re-mapping, all attempts to access the defective row or column will be redirected to the redundant row or column known to be properly working. The remapping process and logical substitution of the redundant region is transparent to the end user. U.S. Pat. No. 6,314,030 to Keeth, the entire disclosure of which is incorporated herein by reference, further describes memory circuit device testing and repair using redundant memory.\nTesting memory devices by writing data to the array is a relatively slow process, especially for flash memory. During such memory testing, one bit pattern which is used employs a repetitive bit pattern, for example, of alternating ones and zeroes arranged in a xe2x80x9ccheckerboardxe2x80x9d or inverse xe2x80x9ccheckerboardxe2x80x9d pattern throughout the memory array. The pattern is written in blocks to the flash memory array.\nThe checkerboard testing pattern is obtained by loading each row of the flash memory array with an alternating series of 1\"\"s and 0\"\"s. The checkerboard pattern alternates from row to row, thus, even rows (rows 0, 2, 4, etc.) are loaded with a repeating pattern of 0\"\"s and 1\"\"s (0-1-0-1 . . . ), while odd rows (1, 3, 5, etc.) are loaded with a repeating pattern of 1\"\"s and 0\"\"s (1-0-1-0 . . . ).\nA short between adjacent cells in two adjacent memory rows is a typical fault in a flash memory array, which requires replacement of the shorted pair of rows with a corresponding pair of redundant rows. Each pair of bad rows will be either an even and odd row (e.g., rows 2 and 3, for example), or an odd and even row (e.g., rows 3 and 4).\nA flash memory device typically is divided into several memory banks, each of which is further divided into a plurality of memory blocks. Each memory block is provided with a number of primary rows and a lesser number of redundant rows, the latter bit available for repairing faulty primary rows in the associated block. Pattern writing during testing includes writing the pattern to the primary rows of the memory block, as well as to the redundant rows of the block, as these too must be tested for integrity. If a defective primary row is replaced by a redundant row during testing, the testing pattern must be rewritten to the block with the redundant row receiving the same bit pattern as the replaced defective row, for further testing. This is a time consuming process.\nThe present invention provides an easier way to test primary rows and redundant rows in memory blocks, and to substitute redundant rows for defective primary rows while permitting continued testing without requiring a rewrite of the testing pattern. In the invention, pairs of defective rows of memory are replaced with redundant rows having an odd-even row pairing which matches the odd-even row pairing of the defective row pair. Consequently, a loaded test pattern of a redundant row used for repair matches that of the defective row which requires repair. This avoids having to rewrite the test pattern when making a repair using redundant rows and saves time in testing the memory.\nA memory device according to the present invention includes memory blocks, each block having rows of primary memory cells having either an even or an odd address. Each memory block also includes at least two rows of redundant memory cells. Like each primary memory row, each row of redundant memory has an even or an odd row address. In each adjacent row pair, whether a primary or redundant row, one row has an even address, and the other row has an odd address.\nDuring testing, the memory block is pre-loaded with an external test pattern, such as a checkerboard pattern, the arrangement of which alternates a xe2x80x9c1xe2x80x9d and xe2x80x9c0xe2x80x9d pattern at each column and with each row, as shown in the example in the Table below:\nUpon detection of a failed odd-even or even-odd pair of primary memory rows, circuitry in the memory device remaps using programmed coding the addresses of the defective pair of adjacent rows, to a corresponding odd-even or even-odd pair of redundant rows as replacements. Advantageously, the replacement redundant rows have addresses that correspond correctly (odd or even) with the odd or even addresses of the replaced primary rows. By using the correct odd-even or even-odd pairing of redundant rows as replacements, the pre-loaded test pattern (checkerboard) will not change and need not be rewritten during further testing, so the pattern that is read out from the replacement redundant rows does not change from that of the replaced primary rows of the block during the further testing. The test pattern used need not be a checkerboard. Any pattern can be utilized that distinguishes even and odd rows, and preferably allows detection of failed row pairs.\nAlthough preferred embodiments of the present invention are described with respect to repairing rows of memory arrays, the invention is not so limited. Repairs of shorted columns, for example, also could take place in odd/even pairs utilizing the teachings of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical tunable filter and a method of manufacturing such an optical tunable filter.\n2. Description of the Related Art\nThere has been known an optical tunable filter for detecting an intensity distribution at each wavelength. The following optical tunable filters have been known in the art.\n1) Filter Formed by Surface Micro-Machining\nIn a conventional optical tunable filter, the thickness of a variable gap is controlled only by the thickness of a sacrifice layer. When the thickness of a variable gap is controlled by the thickness of a sacrifice layer, the variable gap may have variations in thickness depending upon conditions for deposition of the sacrifice layer. Accordingly, uniform Coulomb forces are not produced between a thin film and a drive electrode. As a result, stable driving cannot be achieved.\nFor example, an optical tunable filter disclosed by JP-A No. 2002-174721 has a movable portion projecting on a surface of a substrate. This type of optical tunable filter has a large thickness.\n2) Filter Using SOI Wafer\nFor example, U.S. Pat. No. 6,341,039 discloses that an SiO2 layer of a silicon-on-insulator (SOI) wafer is employed as a sacrifice layer to form a variable gap. According to this method, a variable gap can be formed with accuracy. However, an insulating structure is not provided between a drive electrode and a movable portion. Accordingly, the movable portion and the drive electrode are stuck on each other when a large electrostatic attraction is caused between the movable portion and the drive electrode.\n3) Problems Common to Both Filters\nIn the above filters, a sacrifice layer is finally released to form a variable gap. In this case, a release hole is required to allow a liquid for release to be introduced into the sacrifice layer. Accordingly, an area to which Coulomb forces are applied is reduced to thereby increase a driving voltage. Further, in a case of a small variable gap, sticking is caused when a sacrifice layer is released. Specifically, a thin film and a drive electrode substrate are stuck on each other due to surface tension of a water. Thus, there has been desired a structure that does not need to release a sacrifice layer.\nFurther, because a portion of the movable portion through which light passes is made of silicon, a conventional optical tunable filter cannot be used for purposes other than wavelength separation of infrared light."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a printed circuit board (PCB) for use in chip-on-board (COB) packages and the COB packages using the same. More particularly, the present invention relates to a PCB which has protrusions or trenches within an encapsulation region on the PCB in order to prevent warping of the COB package after the encapsulation process.\n2. Description of the Related Arts\nAs semiconductor integrated circuit devices become more complex and achieve greater speed and functionality, larger numbers of input/output (I/O) pins are required to realize the enhanced performance parameters of the new designs. At the same time, however, the overall size of the devices must continually be reduced to meet the need for employment in smaller and smaller electronic devices. To meet these conflicting needs of higher pin densities in smaller areas, so-called chip-on-board (COB) packages have been developed. In COB packages, a semiconductor integrated circuit device (IC) is directly mounted on a printed circuit board (PCB), thereby eliminating the need to construct a separate lead frame for connection to the device.\nCOB packages are heavily used for memory cards. As is well-known, memory cards have been developed for many uses such as storing captions, sounds, or static images generated by electronic entertainment equipment and digital still cameras, as substitutes for hard disks of computers, and as components of portable information storage and retrieval devices. By packaging several memory chips onto a single board, memory cards can attain a large storage capacity.\nMemory cards include miniature cards, compact flashes, and SmartMedia. Among these, miniature cards and compact flashes suffer from high overhead cost and bulky size because both require built-in controllers.\nIntegrated circuit (IC) cards such as SmartMedia or solid state floppy disc cards (SSFDC), which are used for storing digital signals, contain COB packages with built-in flash memory chips. These IC cards enjoy advantages over earlier memory cards, including smaller size, good portability, and the ability to easily expand or upgrade the memory capacity without having to change the I/O pin name conventions. SmartMedia is used for storing information from, for example, digital still cameras or personal digital assistants (PDAs), and may be used in either an internal or removable mode. Moreover, SmartMedia is expected to be used more extensively in the future for a wide range of applications since it has a greater storage capacity than floppy diskettes or ID cards using a magnetic strip or tape, and is easy to store and handle.\nThe conventional PCB utilized in COB packages is described with reference to a plan view in FIG. 1. The conventional PCB 10 consists of thin board body 11 made from a glass-epoxy resin or BT (Bismaleimide Triazine) resin. A circuit pattern of conductors 13, which serve as electrical signal transmission pathways, is formed in the upper surface of the PCB, for example by using a copper foil and etching process. The bottom surface of the PCB is also provided with a pattern of conducting material which serves as a set of external contacts 15 as shown in FIG. 3. The conductors 13 of the circuit pattern in the upper surface of the PCB 10 are electrically connected to corresponding external contacts 15 by way of via holes 16 in the board body 11.\nThe circuit pattern of conductors 13, the external contacts 15, and the inner wall of the via holes 16 are coated with a conductive plating material, for example gold (Au) which does not suffer oxidation. A chip bonding area 12 is also formed on the upper surface of the board body 11. Usually, a cavity is produced in the bonding area 12 to reduce the thickness of the board body 11 where the chip will be attached and thus reduce the overall height of the finished package.\nThe PCB 10 is usually manufactured in the form of a strip like a conventional lead frame strip. On the strip, each individual PCB 10 is linked to guide rails 19 by way of tie bars 18.\nFIG. 2 and FIG. 3 depict the conventional COB package 100 employing the PCB 10 described above. The COB package 100 includes a semiconductor chip 20, a conventional PCB 10, bonding wires 22 and an encapsulant 23. The chip 20 is attached to the chip bonding area 12 on the upper surface of the PCB 10, and the chip 20 is electrically connected to the circuit pattern of conductors 13 by using bonding wires 22 made from gold (Au) or aluminum (Al). Thus, the chip 20 which is electrically connected to conductors 13 of the circuit pattern via the bonding wires 22, is also electrically connected to the external contacts 15 formed on the lower surface of the PCB by way of the conductors 13 and the via holes 16. The upper surface of the PCB 10 within an encapsulation region is encapsulated with a molding compound, such as a thermosetting compound like epoxy resin, to form an encapsulant 23 which provides protection from external environmental stresses. The encapsulant 23 of the conventional COB package extends above the upper surface of the PCB and is usually formed using a transfer molding process.\nIn an exemplary memory chip configuration, the chip 20 on the chip bonding area 12 may be a flash memory chip having a storage capacity, for example, of 4 Megabytes (M), 16 M or 32 M. In general, the size of a flash memory chip increases with its storage capacity. For example, a 4 M flash memory chip has a size of about 116.1 mm.times.269.1 mm, where mm denotes a millimeter, a 16 M chip has a size of about 196.5 mm.times.444.7 mm, and a 32 M chip has a size of about 294.5 mm.times.514.7 mm. The dimension of the chip bonding area 12 varies according to the size of the chip 20 to be attached to it.\nThe COB package 100 retains the same pin name conventions among the versions using different flash memory chips so that an appliance using the COB package 100 can be easily expanded or upgraded by interchanging cards. Therefore, the structure of the base card on which the COB package 100 will be mounted must be uniform. Consequently, the shape or configuration of the PCB 10 also must be uniform, regardless of the size of the flash memory chip 20. Therefore it is necessary to be able to mount several kinds of flash memory chips, each with different capacity, onto a single sized and shaped PCB.\nSince the shape of the COB package 100 should be uniform regardless of the size of the chip 20 mounted on the PCB 10 therein, the volume of molding compound such as epoxy resin employed for forming the encapsulant 23 of the COB package 100 becomes smaller as the chip size becomes larger. The COB package 100 having the smaller-sized 4 M chip, in which the largest volume of resin is incorporated to form the encapsulant, undergoes more warping than a COB package having the mid-sized 16 M chip. Similarly, the COB package 100 having the larger-sized 32 M chip, in which the smallest volume of resin is incorporated, also undergoes more warping than a COB package having the mid-sized 16 M chip--this time in an opposite direction.\nTherefore, the volume of the molding resin must be balanced with the thermal expansion of the PCB, and with the other conditions of the molding process, in order to prevent failures of the COB package due to warping. However, it is not easy to balance the volumes and prevent such warping failures, because the encapsulant is constrained to have nearly uniform external dimensions while the chips within the encapsulant have different sizes. This situation forces the encapsulant to contain different volumes of molding compound.\nConsequently, the COB package containing a 4 M flash memory chip usually undergoes warping in which the both ends of the PCB bend upward while the center part where the encapsulant is formed bends downward. By contrast, the COB package containing a 32 M flash memory chip undergoes warping in which the both ends of the PCB bend downward while the center part bends upward. Thus warping failures occur more often for COB packages with the smallest or largest memory chips, than for COB packages with a mid-sized memory chip.\nTherefore there is a need for a COB package that reduces warping failures while varying memory chip sizes and simultaneously maintaining nearly uniform encapsulant external dimensions."}
{"text": "For multi-layer IC devices coupled at a frontside to device packaging, e.g., oriented in a flip-chip orientation, debugging for defects in the IC is difficult due to having to approach the desired layers from the backside of the device. FIG. 1 illustrates a sideview block diagram of a typical flip-chip configuration. As shown in FIG. 1, an IC device 10 is coupled to a ceramic package 12 (e.g., a C4 package) via solder bumps 14. The solder bumps 14 act as chip-to-carrier interconnects to attach the IC device 10 to the ceramic package 12 and to mate with corresponding pad patterns to form the necessary electrical contacts between the circuit(s) of the IC device 10 and pins of the package 12.\nTesting of the circuit remains a challenge due to the upside-down nature of the flip-chip orientation. While the circuit may be approached through the backside layers, such techniques are usually not preferred due to the difficulties associated with having to access the layers in an unconventional order. Further, these techniques normally reduce the thickness of the device to reach the circuit, making the device extremely fragile and cumbersome to handle and utilize during testing.\nAccordingly, a need exists for an efficient and effective procedure for separating an IC device from its package to allow utilization, such as testing, from a frontside. The present invention addresses such a need."}
{"text": "1. Technical Field\nThis invention generally relates to a plate washing system with ultrasonic cleaning of pipes and a control method thereof. The ultrasonic cleaning system can provide cleaning of the plate washing system's dispense and aspirate pipes, and monitor the performance in terms of volumetric accuracy and precision.\n2. Discussion of Related Art\nCertain laboratory operations, such as immuno assays, require the testing of small samples which are carried out in an arrangement of microwells having volumes of, for example 50-300 microliters or less formed in microtiter plates, hereinafter referred to generically as well plates. An example of this type of laboratory operation is an Enzyme Linked ImmunoSorbent Assay (“ELISA”) reaction which is performed for measuring the presence or absence of an antigen/antibody complex formed within the wells of the well plate.\nReactions of this type involve the addition and removal of liquid reagents within each well. Intentionally, some of the components in the reagent chemically bond to the well. Therefore, at several stages of the reactions, the unbound liquid and components remaining in the wells must be removed and the insides of the wells must be washed by dispensing a wash solution such as water, a buffer solution, or other fluid in the wells using and gravity feed or a pump, and then evacuating the liquid under a vacuum.\nThe wells can be arranged in a strip or in-line format, or can be arranged in a matrix format. Until recently, commonly used matrices were configured to have 8×12 wells spaced at 9 mm apart between centers, hereinafter referred to as a 96-well plate. However, with the advent of High Throughput Screening (“HTS”), two more matrixes were introduced which increased the total number of wells while keeping the overall size of the well plate the same: 1) the 384-well plate, as shown in FIG. 1, configured to 16×24 wells spaced at 4.5 mm apart between centers, and 2) the 1536-well plate configured to have 32×48 wells spaced at 2.25 mm apart between centers (not shown). Since the overall foot print of the these new well plates are the same as the 96-well plate, the size of the wells in the new micro titer well plates is necessarily smaller than those in the 96-well plates while the depth of the wells remains generally the same. However, this is not always the case.\nA conventional washer used for removing the unbound contents in wells of a well plate includes dispense pipes for dispensing the wash solution into the wells of the plate (e.g., by a pump or gravity feed), and aspirate pipes for evacuation the solution from the wells of the well plate (e.g., by a vacuum or a suction device). In order to quickly wash the well plates, the washing process is performed simultaneously on as many wells of the well plate as possible. A commercial example of such a microplate washer is the Texan PW384.\nAs discussed in U.S. Pat. No. 5,951,783 issued to Kontorovich et al., which is herein incorporated by reference, the dispense and aspirate pipes can be provided on a single manifold assembly or separate dispense and aspirate manifolds.\nIn order to accommodate the well plates having smaller wells, the dispense and aspirate pipes must have small diameters. However, as a result of evaporation, the dispensed materials leave solid materials (such as salts from the assay reagents) within the pipes. The solid material residue can impact the performance of the pipes or even render the pipes inoperable. Impact on performance issues is currently not detectable within the microplate washing system and requires external instrumentation to detect volumetric dispense and aspiration errors.\nUltrasonic cleaning techniques have been used to remove the residual material from the aspirate and dispense pipes and return the pipes to an operative condition. These ultrasonic cleaning techniques use a commercially available ultrasonic tank of suitable size to allow immersion of the impaired pipe assemblies.\nAlthough the use of ultrasonic cleaning is effective, ultrasonic cleaning using a commercially available tank is a complex process. It requires the addition of cleaning liquid in order to fill the tank, removal of the pipe assemblies from the microtiter plate washing system before the pipe assembly is inserted into the tank, removal of the waste material from the tank once the cleaning process has been completed, installation of the pipe assemblies on the plate washing system, and purging cleaning fluid from the pipe assemblies by dispensing."}
{"text": "1. Field of Invention\nThe present invention relates to a high voltage device and a manufacturing method of a high voltage device; particularly, it relates to such device and manufacturing method wherein the breakdown voltage is increased.\n2. Description of Related Art\nFIGS. 1A and 1B show a cross-section view and a 3D (3-dimensional) view of a prior art double diffused drain metal oxide semiconductor (DDDMOS) device, respectively. As shown in FIGS. 1A and 1B, a P-type substrate 11 has multiple isolation regions 12 by which a device region 100 is defined. The isolation region 12 for example is a shallow trench isolation (STI) structure or a local oxidation of silicon (LOCOS) structure, the former being shown in the figures. The DDDMOS device is formed in the device region 100, which includes a gate 13, a drift region 14, a drain 15, and a source 16. The drift region 14, drain 15 and the source 16 are formed by a lithography process and an ion implantation process, wherein the lithography process defines the implantation regions by a photoresist mask together with a self-alignment effect provided by all or part of the gate 13 and the isolation regions 12, and the ion implantation implants N-type impurities to the defined regions in the form of accelerated ions. The drain 15 and the source 16 are beneath the gate 13 and at different sides thereof respectively. Part of the drift region 14 is located beneath the gate 13 near the drain. The DDDMOS device is a high voltage device designed for applications requiring higher operation voltages. However, if it is required for the DDDMOS device to be integrated with a low voltage device in one substrate, the high voltage device and the low voltage device should adopt the same manufacturing process steps with the same ion implantation parameters, and thus the flexibility of the ion implantation parameters for the DDDMOS device is limited; as a result, the DDDMOS device will have a lower breakdown voltage and therefore a limited application range. To increase the breakdown voltage of the DDDMOS device, additional manufacturing process steps are required, that is, an additional lithography process and an additional ion implantation process are required in order to provide different ion implantation parameters, but this increases the cost.\nFIGS. 2A and 2B show a cross-section view and a 3D view of a prior art lateral diffused metal oxide semiconductor (LDMOS) device. Compared to the prior art shown in FIGS. 1A and 1B, the LDMOS device shown in FIGS. 2A and 2B has a body region 17 and a body electrode 18, and part of its gate 13 is located on the isolation region 12. Still similarly, when the LDMOS device is integrated with a low voltage device in one substrate, the high voltage device and the low voltage device should adopt the same manufacturing process steps with the same ion implantation parameters, and thus the flexibility of the ion implantation parameters for the LDMOS device is limited; as a result, the LDMOS device will have a lower breakdown voltage and therefore a limited application range. To increase the breakdown voltage of the LDMOS device, additional manufacturing process steps are required, that is, an additional lithography process and an additional ion implantation process are required in order to provide different ion implantation parameters, but this increases the cost.\nIn view of above, to overcome the drawbacks in the prior art, the present invention proposes a high voltage device and a manufacturing method thereof which provide a higher breakdown voltage so that the high voltage device may have a broader application range, in which additional manufacturing process steps are not required such that the high voltage device can be integrated with and a low voltage device and manufactured by common manufacturing process steps."}
{"text": "The present invention relates to a plasma processing device, and in particular to a plasma processing device using a cathodic arc discharge.\nIn recent years, study of a thin film deposition technique using a low pressure arc discharge has been vigorously conducted. In such a technique, an arc discharge is generated by bringing an electrode generally called striker in mechanical contact with a target portion serving as a cathode or by using an electron beam to let an arc current of several tens amperes flow into the target portion. Plasma is sustained by making ions from plasma hump generated in an upper space of the target collide with the cathode and thereby generating ions and electrons from the cathode. Plasma including these ions and electrons is efficiently led to a vacuum reaction chamber by using a magnetic field duct for transportation and a magnetic field duct for scanning. The plasma is thus applied uniformly to a substrate to be processed. Thereby, processing such as thin film forming or etching is conducted.\nWhen plasma is generated by an arc discharge in the above described conventional method, however, a large number of neutral particles having no electric charge or electrically charged particles are generated besides ions and electrons. These hamper the processing such as thin film forming or etching as particles, resulting in a great problem."}
{"text": "One or more embodiments relate to the technology of fabricating a semiconductor device, and more particularly, to a semiconductor device with reduced contact resistance between a substrate and a plug, and a method of fabricating the semiconductor device.\nAs the degree of integration of semiconductor devices increases, a space between gate patterns becomes narrow and thus a contact processing margin reduces. Therefore, in order to secure a desired contact processing margin, a landing plug contact (LPC) structure is commonly used.\nA landing plug contact process is accomplished by filling a space between gate patterns where, previously, a bit line contact and a storage contact are formed with a conductive layer to secure an overlay margin in subsequent contact processes.\nFIG. 1 illustrates a cross-sectional view of a conventional semiconductor device.\nReferring to FIG. 1, a gate 17 is formed by sequentially stacking a gate insulation layer 14, a gate electrode 15 and a gate hard mask layer 16 over a substrate 11 including a device isolation region 12 and an active region 13, and then etching the gate insulation layer 14, the gate electrode 15 and the gate hard mask layer 16 at once by using a gate mask. Then, after forming gate spacers 18 on sidewalls of the gate 17, a source and drain region (not shown) is formed in the active region 13 between two neighboring gates 17.\nAfter depositing an interlayer insulation layer (not shown) covering the resultant structure, the deposited interlayer insulation layer undergoes a photolithography process and an etching process using a contact mask as an etch barrier. As a result, a contact hole is formed. The contact hole is then filled with a conductive layer, thereby forming a landing plug 19.\nHowever, during the formation of the gate 17, since etch rates of the gate insulation layer 14, the gate electrode 15 and the gate hard mask layer 16 with respect to an etching gas are different from each other, the resultant gate 17 is formed having a sidewall with a positive slope. That is, a bottom line width (bottom CD) W2 of the gate 17 is greater than its top CD W1 (W1<W2). Therefore, a space between two neighboring gates 17 is reduced and thus a contact area between the landing plug 19 and the substrate 11 is also reduced. The reduction of the contact area between the landing plug 19 and the substrate 11 becomes a more serious issue as the gate spacers 18 are formed. Byproduces or residues R generated during formation of the gate 17 may further reduce the contact area between the landing plug 19 and the substrate 11.\nAs described above, when the contact area between the landing plug 19 and the substrate 11 is reduced, the contact resistance between the landing plug 19 and the substrate 11 is increased, thus, the operational speed of the semiconductor device deteriorates."}
{"text": "It is often desirable to mount point-of-purchase display devices on the modular shelving that is used in retail stores. The commonest types of such shelving have hollow steel upright supporting posts at each end of each module, and each supporting post has a line of spaced longitudinal slots in which mounting hooks for the merchandise supporting shelves are engaged. In most cases, the posts of laterally adjacent modules abut one another; although there are certain types of shelving in which a single post supports two adjacent modules.\nVarious manufacturers of modular shelving use supporting posts of different cross sectional configurations and areas, and the slots in the posts of different manufacturers' shelving are of different lengths and different spacings. Apparatus for mounting point-of-purchase display devices on such modular shelving must be capable of being mounted on the supporting posts of any of the several types of shelving used in retail stores, in order that the company placing the display devices need not have different mounting brackets for each different kind of shelving.\nThere is one company which markets a \"universal\" mounting key which is sufficiently universal that only two types of such mounting keys are required to fit many of the slotted hollow steel posts now on the market. Each mounting key is fabricated from a piece of 3/32\" (2.38 mm) sheet metal which may be inserted edgewise through a post slot. The key includes a shank which is threaded top and bottom to receive a wing nut, and at one end of the shank are ears which engage the interior of a post above and below a post slot through which the key is inserted. The wing nut may be screwed into firm clamping engagement with the post wall, or with a bracket which is supported upon the shank and clamped against the post wall by the wing nut.\nThe above-described mounting keys have proved to be satisfactory for connecting supporting brackets to the upright hollow posts of retail shelving. However, the wing nuts which are used on such mounting keys must be of fairly large diameter to function properly, and the span across the wings must be large enough that two such keys may not be used in coplanar slots of the laterally abutting posts of two adjacent modules because the wind nuts interfere with one another; nor, of course, may they be used in a single slot of a post that supports two modules.\nThere are certain situations in which it is desirable to have several point-of-purchase display devices placed end-to-end above adjacent modules, and there is no way that the above-described mounting keys can be used for this purpose."}
{"text": "1. Field of the Invention\nThis invention relates to an improvement in a camera of the kind using a film provided with a magnetic storage part and having a magnetic head arranged in the camera to read information such as ISO data, etc. written in the magnetic storage part of the film and to write therein some information as necessary.\n2. Description of the Related Art\nCameras of the kind using a kind of film that has a magnetic storage part have recently been proposed, for example, as disclosed in U.S. Pat. No. 4,864,332. This camera is arranged to record in the magnetic storage part some photographic information such as a shutter time value, an aperture value, a date, a title, etc. and to read the recorded information as necessary out of the magnetic storage part by means of a magnetic head. However, practical attempts for affecting cameras of the kind having this magnetic head have encountered the following problems:\n1) Cameras of the kind disclosed in U.S. Pat. No. 4,841,319 etc. are arranged as follows: A film cartridge containing the whole roll of film completely therein including the leader part of the film is used. The film cartridge is inserted into a cartridge chamber, for example, from one side of the camera. The film is then pushed out of the cartridge by driving its shaft with a fork. With the cartridge of the kind used for a film having a magnetic storage part, an automatic film loading operation tends to become impossible, because the writing/reading magnetic head is merely arranged to come into contact with the film surface, the fore end part of the film tends to collide either with the magnetic head or with a film retaining pad and to be bent and warped when it is pushed out from the cartridge.\n2) When the camera is carried by the photographer, for example, by hanging on the neck while walking, vibrations are transmitted to the camera. Under such a condition, various accidents tend to happen including that the magnetic storage part of the film might be damaged by the movement of the magnetic head."}
{"text": "The present invention relates to drill pipe and more particularly to multiple wall drill pipe having means for releasably connecting an inner pipe member to an outer pipe member of the multiple wall drill pipe.\nMultiple wall drill pipe, e.g., such as that shown in U.S. Pat. Nos. 1,547,461; 1,909,075; 2,850,264; 3,208,539; and 3,664,441, has been used for many years in both conventional and reverse circulation drilling operations. Such multiple wall pipe is conventionally made of an outer pipe member and an inner pipe member disposed concentrically within the outer pipe member and longitudinally fixed thereto such as by straps welded to the inner wall of the outer pipe member and the outer wall of the inner pipe member at one or more locations along the length of the multiple wall pipe. Sections of the multiple wall pipe are adapted to be joined in longitudinal alignment to form a drill pipe string extending from a cutting bit to the surface of the ground where the drilling operation is being conducted. Two continuous flow passages are formed in the multiple wall pipe, a central passage defined by the inner surfaces of the inner pipe members and an annular passage formed between the outer surfaces of the inner pipe members and the inner surfaces of the outer pipe members. These passages are used to transport a drilling fluid from the ground surface to the drill bit and then back to the surface. For example, in a reverse circulation drilling operation, drilling fluid is injected into the upper end of the annular passage under pressure, is forced downwardly therethrough to the drill bit and then returns to the surface through the central passage along with fragments generated by the drilling operation.\nIn rotary drilling operations using multiple wall drill pipe, torque is transmitted to the drill bit through the outer pipe members. It has been determined that if the inner pipe members are fixedly non-rotatably rigidly connected by welds to the outer pipe members, relatively high torsional stresses may be transmitted from the outer pipe member to the inner pipe member through the fixed non-rotatable rigid connections during rotary drilling operations with resulting pipe failures caused by torsional stress and metal fatigue. Such failures usually occur in an area at or adjacent the fixed welded connection between the inner and outer pipe members. In some situations, such failures may be repaired by replacement of the inner pipe member and rewelding of the fixed conditions. However, to the extent that repair of damaged sections of conventional multiple wall pipe has been possible at all, such repair has been extremely difficult and costly. Typical procedures for repair of broken or damaged connections between the inner and outer pipe have involved wall cutting operations such as with a torch, to effect necessary repair or replacement of component parts. This procedure is frequently found to be either mechanically or economically prohibitive, resulting in complete replacement of the damaged section with a new section of multiple wall pipe without salvage of either the inner or outer pipe member.\nThe present invention involves new and improved connecting means for releasably and rotatably connecting the inner pipe member to the outer pipe member to enable repair and/or replacement while holding the inner pipe member in fixed axial relationship under force loads applied in use during a drilling operation and while enabling relative rotary movement between the inner pipe member and the outer pipe member to relieve torsional stresses encountered during use, the releasable connecting means permitting the removal of the inner pipe member from the outer pipe member by relative axial displacement therebetween caused by application of only an axially directed force load in excess of the axially directed force loads applied to the pipe during use in a drilling operation. In general, the connecting means comprises a plurality of axially spaced sets of resilient compressible axially spaced connecting members mounted in corresponding axially spaced sets of axially spaced inner and outer groove means providing a plurality of axially interengaged stop surfaces which prevent relative axial movement between the inner pipe and the outer pipe in either direction occasioned by any axial forces applied thereto in use during a drilling operation whereby the inner pipe is located in fixed axial position relative to the outer pipe at more than one location along the length of the pipes.\nIn the illustrative and presently preferred embodiment of the invention there are at least two sets of connecting members and associated grooves with one set located adjacent one end of the pipe section and the other set located adjacent the other end of the pipe section to axially fix each end portion of the inner pipe member relative to each end portion of the outer pipe member. Each set of connecting members and associated grooves comprise three equally axially spaced inner groove means formed in the outer periphery of the inner pipe and three similarly equally axially spaced outer groove means formed in the inner periphery of a connecting sleeve means fixedly attached to the outer pipe member. The connecting members, which are resilient compressible O-ring members, are mounted in corresponding axially aligned ones of the inner and outer connecting groove means during assembly so as to connect and support the inner pipe member relative to the outer pipe member after assembly and during use in a manner preventing relative axial displacement while permitting relative rotational displacement during use and while enabling disassembly of the inner and outer pipe for repairs or the like upon application of axially directed forces in excess of any axially directed forces normally encountered during usage."}
{"text": "The disclosure relates to the field of biosensors for label independent detection (LID). More particularly the disclosure relates to optically scanned label independent detection biosensors and to a system and method for microplate image analysis."}
{"text": "U.S. Pat. Nos. 5,100,354, 5,249,998 and 6,179,676 B1 disclose a water sport device for supporting a seated human rider while the rider and the device are towed behind a powered water craft. The device includes an elongated board to which a rigidly mounted seat and foot holders are secured. An elongate strut in the shape of a hydrofoil projects downwardly from the board. The lower end of the strut carries an arm which is generally parallel to the board. Front and rear planing blades are carried by the arm. The positioning of the rigidly mounted seat and the planing blades, the use of a single vertical strut, the size of the planing blade and the positioning of the foot holders provides a water sports device which is relatively easy to ride, while at the same time being highly maneuverable and capable of high jumps.\nThe present invention provides significant improvements in this water sport device.\nFirst Embodiment\nIn a water sports device for supporting a seated human rider while the rider and the device are towed behind a powered water craft, the device including an elongated board having a front end and a back end to which is secured a seat for supporting the buttocks of the rider in a position spaced from and roughly centered above the back one-third of the board, a holder for securing the feet of the rider over the top of the board secured to the board spaced at least about two feet toward the front end of the board from the seat, an elongated hydrofoil extends downward from the board, front and rear planing blades carried by an arm which is generally parallel to the board, said arm being secured to the hydrofoil and spaced from the board, so that the blades are generally parallel to the board, the planing blades provide essentially no lift when the board is horizontal;\nthe improvement wherein the distance between the front and rear planing blades can be changed fore and aft by adjustment of the position of the rear blade so that the advanced rider can move the rear blade forward to make the device more responsive and thereby create higher jumps, and the less experienced rider can move the rear blade further to the rear of the front blade to make the device less responsive.\nSecond Embodiment\nIn a water sports device for supporting a seated human rider while the rider and the device are towed behind a powered water craft, the device including an elongated board having a front end and a back end to which is secured a seat for supporting the buttocks of the rider in a position spaced from and roughly centered above the back one-third of the board, a holder for securing the feet of the rider over the top of the board secured to the board spaced at least about two feet toward the front end of the board from the seat, an elongated hydrofoil extends downward from the board, front and rear planing blades carried by an arm which is generally parallel to the board, said arm being secured to the hydrofoil and spaced from the board, so that the blades are generally parallel to the board, the planing blades provide essentially no lift when the board is horizontal;\nthe improvement wherein said rear planing blade has an angle which slopes downward from the top surface to its leading edge whereby the device is more stable and less prone to porpoising.\nThe first improvement by creating an adjustable rear blade enables an advanced rider to move the rear blade forward on the arm to reduce the space between the planing blades. This makes the entire water sports device more responsive so that when the rider leans back, the device comes out of the water and the height of the jump is greater. Conversely, for the less experienced rider, the spacing between the blades can be increased, which increases the stability by reducing the responsiveness.\nThe second improvement pertains to an improvement in the shape of the rear blade. In the prior water sports device, the rear blade sloped upward from the underside of the blade to the leading edge. In the present invention, the angle is downward at an increasing angle on the top surface of the rear planing blade to the leading edge. The trailing edge is thinned out on the underside of the rear planing blade. The new downward angle stabilizes the rider by creating a torque or xe2x80x9ctow inxe2x80x9d that reduces the porpoising or up and down pitching effect observed in riding the devices of the prior art referred to hereinabove."}
{"text": "Referring to FIG. 1, the nasal airway 1 comprises the two nasal cavities separated by the nasal septum, which airway 1 includes numerous ostia, such as the paranasal sinus ostia 3 and the tubal ostia 5, olfactory cells and is lined by the nasal mucosa. The nasal airway 1 can communicate with the nasopharynx 7, the oral cavity 9 and the lower airway 11, with the nasal airway 1 being in selective communication with the anterior region of the nasopharynx 7 and the oral cavity 9 by opening and closing of the oropharyngeal velum 13. The velum 13, which is often referred to as the soft palate, is illustrated in solid line in the closed position, as achieved by providing a certain positive pressure in the oral cavity 9, such as achieved on exhalation through the oral cavity 9, and in dashed line in the open position.\nThere are many nasal conditions which require treatment. One such condition is nasal inflammation, specifically rhinitis, which can be allergic or non-allergic and is often associated with infection and prevents normal nasal function. By way of example, allergic and non-allergic inflammation of the nasal airway can typically effect between 10 and 20% of the population, with nasal congestion of the erectile tissues of the nasal concha, lacrimation, secretion of watery mucus, sneezing and itching being the most common symptoms. As will be understood, nasal congestion impedes nasal breathing and promotes oral breathing, leading to snoring and sleep disturbance. Worryingly, the incidence of such allergic and non-allergic inflammatory diseases is increasing. Other nasal conditions include nasal polyps which arise from the paranasal sinuses, hypertrophic adenoids, secretory otitis media, sinus disease and reduced olfaction.\nIn the treatment of certain nasal conditions, the topical administration of medicaments is preferable, particularly where the nasal mucosa is the prime pathological pathway, such as in treating or relieving nasal congestion. Indeed, topical administration is advantageous in minimizing the possible side effects of systemic administration. Medicaments that are commonly topically delivered include decongestants, antihistamines, cromoglycates, steroids and antibiotics.\nThere are now an increasing number of adults and children who rely on pharmaceuticals to relieve symptoms associated with nasal conditions. At present, among the known anti-inflammatory pharmaceuticals, topical steroids have been shown to have an effect on nasal congestion. Topical decongestants have also been suggested for use in relieving nasal congestion. The treatment of hypertrophic adenoids and chronic secretory otitis media using topical decongestants, steroids and anti-microbial agents, although somewhat controversial, has also been proposed. Further, the topical administration of pharmaceuticals has been used to treat or at least relieve symptoms of inflammation in the anterior region of the nasopharynx, the paranasal sinuses and the auditory tubes.\nAside from the delivery of medicaments, the irrigation of the nasal mucosa with liquids, in particular saline solutions, is commonly practiced to remove particles and secretions, as well as to improve the mucociliary activity of the nasal mucosa. These solutions can be used in combination with active pharmaceuticals.\nFurthermore, medicaments are now increasingly systemically delivered through the nasal pathway, the nasal pathway offering a good administration route for the systemic delivery of pharmaceuticals, such as hormones, for example oxytocin, and antimigraine compositions, as the high blood flow and large surface area of the nasal mucosa advantageously provides for rapid systemic uptake.\nA variety of delivery systems have been developed to deliver substances to the nasal airways of subjects.\nConventionally, spray bottles have been used to deliver a medicament-containing liquid or an irrigating liquid to the nasal airways of subjects. However, the distribution of the delivered substance, in particular to the posterior region of the nasal airway, is less than ideal, especially in the cases of moderate and severe nasal obstruction. This poor distribution is often further exacerbated by a subject inhaling through the nasal airway during delivery, as is often prescribed, in an attempt to deliver the substance to the posterior region of the nasal airway. Indeed, an amount of the substance can be drawn into the lungs or swallowed in each delivery, which could be problematic in paediatric subjects if the medicament is a potent pharmaceutical, such as a steroid, which has to be administered frequently. In addition, the spray is frequently directed against the nasal septum which can undesirably lead to localized deposition. Further, the mechanical action of the delivery mechanism of the spray bottles can cause irritation and bleeding.\nGB-A-408856 discloses a delivery device, which, in one mode of use, allows for the inhalation of separate air flows entraining medicament into respective ones of the nasal cavities of a subject. This delivery device comprises a chamber containing a sponge saturated with medicament, a mouthpiece connected to the chamber and first and second nosepieces connected to the chamber. In one mode of use, the nosepieces are fitted into respective ones of the nostrils of a subject, and, on inhalation through the nosepieces, air flows entraining medicament are drawn into the lungs of the subject. In another mode of use, the mouthpiece is taken in the mouth of a subject, and, on inhalation through the mouthpiece, an air flow entraining medicament is drawn into the lungs of the subject.\nWO-A-98/53869 discloses a delivery device for delivering a powder containing a medicament to the nasal mucosa in one of the nasal cavities of a subject. This device comprises a tubular section which contains a metered dose of powdered medicament. In use, the ends of the tubular section are respectively located in the nostril of one of the nasal cavities and the mouth of a subject, and on exhalation by the subject through his or her mouth the exhaled air entrains the powdered medicament and delivers the same into the one nasal cavity, with the exhaled air backflowing out of the one nostril around the tubular section. In one embodiment the tubular section includes a flexible portion upstream of the dose of powdered medicament. The provision of this flexible portion allows the subject to close the tubular section at a point upstream of the medicament, such that, on release of the closed flexible portion during exhalation, a short explosive air flow entraining medicament is delivered into the one nasal cavity. In another embodiment the end of the tubular section located in the nostril can be shaped to act to locate the tubular section in a position in the nostril which allows for the deposition of the powdered medicament on the nasal mucosa.\nWhilst this delivery device is simple in construction, the operation of the device still does not provide for the effective delivery of substances, in particular one of a liquid or powder containing medicament, to the posterior region of the nasal airway, since medicament is delivered separately to each of the nasal cavities and the air flow into and out of each nasal cavity is through the same opening, namely the respective nostril, with the closed posterior region of the respective nasal cavity acting as a pressure reflecting surface which causes the exhaled air to backflow out of the one nostril before ever adequately reaching the posterior region of the respective nasal cavity. Further, in providing a short explosive burst of air flow into one of the nasal cavities, it is not possible to achieve a sustained and controlled bidirectional air flow through the nasal cavities which has been found necessary to deliver a substance effectively to the posterior region of the nasal airway.\nFor any substance to be delivered effectively to the nasal airway, it is highly desirable that the administration is efficient and simple. However, there can be problems in attempting to achieve this goal. In particular, the pathological changes observed with nasal inflammation make administration of substances, such as liquids or powders, tricky, particularly to the posterior region of the nasal airway and the posterior margins of the nasal structures. Indeed, as a consequence of the complex geometry of the narrow slit-like passages in the nasal airway, these passages become partially occluded when the nasal mucosa is inflamed and congested, making the distribution of topical pharmaceuticals to the nasal airway difficult."}
{"text": "A tape feeder is frequently used as a component supply device in a component mounting device. The tape feeder feeds a carrier tape which holds components from a tape reel and pitch-feeds the carrier tape in synchronism with timing for mounting the component. As a reel exchange operation performed when component depletion occurs while the tape feeder is performing mounting operation, a splicing method for splicing carrier tapes with each other has been adopted. In the method, a tail end of a carrier tape already loaded and a head of a new carrier tape are spliced together. Consequently, leading end placing operation for placing the leading end of the new tape on the tape feeder and positioning the tape can thereby be omitted. Accordingly, it is possible to omit a wasteful time caused by stopping a mounting device every time components become depleted. Further, a splice between two carrier tapes is automatically detected at a specific position on the tape feeder, thereby providing the following advantage. That is, automatic detection of the splice can be applied to managing the number of remaining components and managing a component supply history in the component supply, such as prediction of timing at which operation for splicing a current carrier tape to the next carrier tape becomes necessary (see, for example, Patent Document 1).\nSpecifically, in a component supply system using tape feeders, component information about a carrier tape unreeled from an individual tape reel and mounted on a corresponding tape feeder is stored in a component information storage unit in real time. The component information includes, for example, a reel ID code for specifying the carrier tape, the number of remaining components not picked up yet, etc. Data rewriting processing for the stored contents on detection of a splice in the carrier tape as a trigger signal. The data rewriting processing includes, for example, operation for writing data by a reel ID code which specifies a carrier tape to be newly spliced, operation for resetting the number of remaining components, etc. Accordingly, it is possible to automatically perform highly reliable component supply management without inconvenience to a machine operator."}
{"text": "1. Field of the Invention\nThe present invention relates to a nitride semiconductor light emitting device; and, more particularly, to a nitride semiconductor light emitting device capable of improving an ESD (Electrostatic Discharge) tolerance (withstand voltage characteristic) of a GaN based nitride light emitting device and a method for manufacturing the same.\n2. Description of the Related Art\nGenerally, it has been known that a GaN based light emitting device has a bad static electricity characteristic in comparison with other compound light emitting devices. This reason is that a lot of crystal defects are formed in a GaN thin film due to a great lattice mismatch (16%) between a substrate and a thin film to be grown since the GaN light emitting device is formed on a sapphire substrate having a great lattice mismatch.\nThis crystal defect increases a leakage current of device and an active layer of the light emitting device having a lot of crystal defects is broken by a strong field when an external static electricity is applied. Generally, it has been known that the crystal defect (threading dislocations) in the order of 1010˜1012/cm2 exist at the GaN thin film.\nThe static electricity broken characteristic of the light emitting device is a very important issue related to the application range of the GaN based light emitting device. Specifically, a design of devices for withstanding a static electricity generated from package devices and workers of the light emitting device is a very important parameter for improving yield of a final device.\nParticularly, the static electricity characteristic has been become more important since the GaN based light emitting device are recently used by being applied to a bad condition in environment such as an outdoor signboard and a vehicle light.\nGenerally, an ESD (Electrostatic Discharge) of a conventional GaN light emitting device withstands several thousands volts in a forward direction under HBM (Human Body Mode), whereas it does not withstand several hundreds volts in a reverse direction. The reason is that the crystal defects of the device are major reasons as described above, and also, an electrode design of the device is very important. Particularly, since the GaN light emitting device generally employs a sapphire substrate as an insulator, the ESD characteristic is further deteriorated by increasing a concentration phenomenon of currents around an N-electrode during a practical device operation as the N-electrode and a P-electrode are formed at the same plane on the structure of the device.\nA conventional method for improving such ESD characteristic has been frequently approached in external aspects of the device. U.S. Pat. No. 6,593,597 teaches a technique for protecting a light emitting device from the ESD by connecting an LED and a Schottky diode in parallel by integrating the LED device and the Schottky diode on the same substrate. In addition, in order to improve the ESD tolerance, a method for connecting an LED to a Zenor diode in parallel has been proposed. However, since such methods cause complicatedness to be assembled by purchasing an additional Zenor diode or performing a Schottky junction and increase the manufacturing cost of the device, they are not preferable in point of cost or yield.\nThe most preferable method is that the ESD characteristic is improved by the light emitting device in itself by improving a thin film characteristic or a structure of the light emitting device. Japanese Pat. No. 3,622,562 discloses a method for improving an ESD characteristic by using a multi-stacked thin film structure using un-doped nitride indium gallium and un-doped nitride gallium on an n-type contact layer, whereas Japanese Pat. No. 3,063,756, Japanese Pat. No. 3,424,629, Japanese Pat. No. 2006-237254 and U.S. Pat. No. 6,677,619 try to reduce leakage currents and improve ESD tolerance by making a layer structure a good crystal property by reducing the density and size of defects to be generated from each layer of the light emitting device.\nAs like this, it is preferable that the quality of a GaN thin film is intrinsically increased to improve the ESD characteristic, but there exists a limitation to remove the defects, therefore, a method of growing the thin film has been continuously continued to improve the ESD characteristic of the GaN up to date."}
{"text": "1. Field of the Invention\nThe present invention relates to an ink jet recording sheet. More particularly, the present invention relates to an ink jet recording sheet which has excellent ink jet recording properties and has excellent gloss of the ink image-formed portions and the non-image-formed portions thereof.\n2. Description of the Related Art\nThe recording by an ink jet printer is often utilized in many fields because printing noise is low, the recording can be affected at a high speed and full-colored images can be easily formed. As recording sheets for the ink jet printing system, wood-free paper sheets modified so as to enhance the ink absorption thereof, and coated paper sheets having a coating surface layer containing porous pigment particles are used. These conventional recording paper sheets mostly have a low surface gloss and a matted appearance. Currently, new types of ink jet recording sheets having a high surface gloss and good appearance are in demand.\nGenerally, as printing paper sheets having a high surface gloss, coated paper sheets having a high gloss and produced by coating a surface of a paper sheet substrate with a coating layer containing plate-crystalline pigment particles and optionally applying a calendering treatment to the coating layer and cast-coated paper sheet produced by coating a surface of a paper sheet substrate with a coating liquid, pressing the coating liquid layer against a mirror-finished surface of a heating drum and drying the coating liquid layer on the mirror-finished surface of the heating drum so as to transfer the mirror surface to the resultant coating layer, are used.\nThe cast-coated paper sheets are mostly used for high gloss printing, because they have higher surface gloss and surface smoothness than those of the usual calender-finished coated paper sheets, and thus render excellent printing effect. However, when the cast-coated paper sheets are used for the ink jet recording system, various difficulties occur.\nA conventional cast-coated paper sheet is disclosed in U.S. Pat. No. 5,275,846. In this cast-coated paper sheet, a high gloss is obtained by transferring the mirror-finished drum surface of a cast-coater to a coating layer containing a pigment and a film-forming binder. This film-forming binder causes the resultant coating layer to be non-porous and thus to exhibit a significantly reduced ink absorption when the cast-coated paper sheet is used for the ink jet recording system. To enhance the ink absorption, it is necessary to make the coating layer porous, so that the resultant porous coating layer can easily absorb the ink. To enhance the porosity of the coating layer, it is necessary to reduce the film-forming property of the coating layer. When the content of the film-forming binder in the coating layer is reduced, the resultant cast-coated paper sheet exhibits a reduced gloss. Therefore, it is extremely difficult to obtain a cast-coated paper sheet having both a high gloss and a satisfactory applicability to the ink jet recording.\nTo solve the above-mentioned problem, Japanese Unexamined Patent Publication No. 7-89,220 discloses a cast-coated paper sheet, having both excellent gloss and superior ink absorption, usable for an ink jet recording system. The cast coated paper sheet is produced by undercoating a surface of a paper sheet substrate with a undercoat layer comprising, as principal components, a pigment and a binder, coating the undercoat layer surface with a coating liquid containing, as a principal component, a copolymeric resin produced by polymerizing at least two monomers having an ethylenically unsaturated bond and having a gloss-transition temperature of 40xc2x0 C. or more, to form a cast-coating liquid layer, pressing the cast-coating liquid layer in wet condition, against a mirror-finished casting drum surface and drying the cast coating liquid layer on the casting drum surface to form a cast-coated layer.\nCurrently, due to the rapid speed of ink jet printers, the high precision of the printed images and the full colored images, the use of the ink jet printing system is spreading, and thus the ink jet recording sheets are required to have a high gloss and improved qualities of the received ink images, namely a high color density and clarity. For example, high gloss and image quality comparative with those of silver salt type photographic prints are required from the ink jet recording sheets. To meet with the requirements, the gloss and printing performance of the ink jet recording sheets must be further improved.\nJapanese Unexamined Patent Publication No. 7-329,412 discloses an ink jet recording sheet in which a cationic compound is contained in a substrate and an ink-receiving layer. In this case, the cationic compound is used for the purpose of preventing migration or permeation of ink in the printed portions of the recording sheets. This publication is quite silent as to the gloss of the ink jet recording sheet.\nAn object of the present invention is to provide an ink jet recording sheet having excellent gloss of ink-image-formed portions and non image-formed portions, superior ink absorption and a satisfactory water resistance of images and capable of recording ink images with a high color density and clarity.\nThe above-mentioned object can be attained by the ink jet recording sheet of the present invention which comprises\na substrate sheet;\nat least one ink receiving layer formed on at least one surface of the substrate sheet and comprising a pigment and a binder; and\nat least one gloss layer formed on the ink receiving layer and comprising a polymer resin,\nat least one cationic compound being impregnated in the substrate sheet or coated on at least one surface of the substrate sheet on which the ink receiving layer is formed.\nThe ink jet recording sheet comprises a substrate sheet, at least one ink receiving layer formed on at least one surface of the substrate sheet and at least one gloss layer formed on the ink receiving layer.\nThere is no limitation to the type and dimensions of the substrate sheet. Usually, the substrate sheet is selected from paper sheets including acid paper sheets and neutralized paper sheets, synthetic paper sheets and air-permeable polymer films. Especially, the paper sheets are preferred for the substrate sheet.\nIn the ink jet recording sheet of the present invention, a cationic compound is impregnated in the substrate sheet or coated on a surface of the substrate sheet on which surface the ink receiving layer is arranged. The cationic compound usable for the present invention is not limited to a specific class of compounds as long as the compound molecule has a cationic moiety. To cause the cationic compound to be retained in or on the substrate sheet with a high efficiency, however, the cationic compound is preferably selected from cationic oligomers and polymers having a high molecular weight of, for example, 1000 or more. There is no specific upper limit to the molecular weight. To stop the cationic compound-containing coating solution having too high a viscosity, the upper limit of the molecular weight is preferably 1,000,000 or less.\nThe cationic compounds usable for the present invention include cationic resins (cationic polymeric compounds) and low molecular weight cationic compounds (cationic monomeric compounds), for example, cationic surface active compounds. To retain the cationic compounds on the substrate sheet such as paper sheet with a high efficiency, the cationic compounds are preferably selected from the cationic resins which may be practically used in the form of an aqueous solution or an aqueous emulsion. There is no limitation to the molecular weight of the cationic resins. In view of the ease of handling, coating and impregnating, however, the molecular weight of the cationic resins is preferably about 1000 to about 1,000,000.\nAlso, the cationic resin may be used in the form of cationic organic pigment particles which have been prepared by insolubilizing the cationic resins by, for example, reacting them with a cross-linking agent into the form of fine particles. The cationic pigment particles can be produced by copolymerizing a cationic comonomer with a polyfunctional comonomer to form a cross-linked polymer particle, or by cross-linking a cationic resin having cross-linking reactive groups, for example, hydroxyl, carboxyl, amino, and/or acetacetyl groups and optionally a cross-linking agent, by means of heating or radiation, to provide a cross-linked resin.\nThe cationic resins usable for the present invention include:\n1) polyalkylenepolyamines, for example, polyethylenepolyamines and polypropylenepolyamines, and derivatives thereof;\n2) acrylic resins having at least one member selected from secondary amine groups, tertiary amino groups and quaternary ammonium groups;\n3) polyvinylamines and polyvinylamidines;\n4) dicyandiamide resins, for example, dicyandiamide-formaldehyde polycondensation products, dicyandiamidediethylenetriamine polycondensation products and polyalkylenepolyaminedicyandiamide copolymers;\n5) epichlorohydrin-dimethylamine addition polymerization products;\n6) dimethyldiallyl ammonium chloride-SO2 copolymerization products;\n7) diallylamine-SO2 copolymerization products;\n8) dimethyldiallyl ammonium chloride polymerization products;\n9) allylamine salt polymerization products;\n10) dialkylaminoethyl(meth)acrylate quaternary salts polymerization products;\n11) acrylamide-diallylamine salt copolymerization products; and\n12) polyaminoepoxy resins, ring-opening polymerization products of ethyleneimine, homo- and co-polymers of cationic vinyl monomers and chitosan.\nAmong the above-mentioned cationic compounds, the dicyandiamide resins are preferably used for the present invention, because they exhibit excellent effects desired for the present invention.\nEspecially, the copolymers of polyalkylenepolyamines with dicyandiamide are preferred. For example, cationic resins consisting of copolymers of at least one member selected from straight chain polyamines, for example, diethylenetriamine, triethylenetetramine, tetraethylenepentamine and iminobispropylamine and organic and inorganic salts, for example, hydrochloric acid salts, sulfates and acetates thereof, with dicyandiamide, are more preferred.\nIn the application of the cationic compound to the substrate sheet, a solution of the cationic compound in water or alcohols is impregnated in or coated on a surface of the substrate sheet by using a size-press machine, gate roll coater, blade coater, air knife coater, roll coater, bar coater or gravure coater. The impregnation or coating of the cationic compound solution may be affected within a paper-forming system (paper-making machine) or outside of the paper-forming system. When the impregnation or coating is carried out, one or both of the surfaces of the substrate sheet are coated with the cationic compound. Also, the cationic compound is applied to at least one surface of the substrate sheet on which surface the ink receiving layer and then the gloss layer are formed.\nUsually, the cationic compound is applied preferably in a dry amount of 0.1 to 10 g/m2, more preferably 0.2 to 5 g/m2 to each of the substrate sheet. If the cationic compound amount is less than 0.1 g/m2, the desired effect may not be obtained. Also, an amount more than 10 g/m2 may cause the desired effect to be saturated and an economical disadvantage may occur. The cationic compound may be applied alone or optionally in a mixture with a binder to the substrate sheet. Optionally, a resin which is usually used as a binder is mixed to the coating liquid or impregnating liquid containing the cationic component. The mixed binder contributes in adjusting the viscosity and coverage of the cationic compound coating liquid and to enhancing the water resistance of the substrate sheet, and surface strength and interlaminar strength of the recording sheet. Also, in the drying step of the size-press procedure the mixed binder enables the cohesion of the recording sheet to the cylinder dryer to be prevented and the resultant recording sheets to exhibit a uniform quality even when the coating amount is fluctuated.\nThe above-mentioned binder contains at least one member selected from, for example, proteins, for example, casein, soybean protein and synthetic proteins, starch compounds, for example, starch and oxidized starches, polyvinyl alcohols, cellulose compounds, for example, carboxymethylcellulose and methylcellulose, conjugated diene polymer latices, for example, styrene-butadiene copolymer latices and methylmethacrylate-butadiene copolymer latices, acrylic polymer latices, and vinyl polymer latices, for example, ethylene-vinyl acetate copolymer latices, which are usually used as coating agent for paper sheets. These binders can be used alone or in combination of two or more thereof.\nWhen the cationic compound is used together with the binder, there are no limitations to the types and coverages thereof. Usually, the cationic compound and the binder are employed preferably in a mixing ratio in weight of 10/90 to 90/10, more preferably 20/80 to 80/20. Optionally, the cationic compound is employed in combination with a pigment, for example, usable for the ink receiving layer.\nWhen the substrate sheet consists of a paper sheet, degree of sizing of the cationic compound-applied paper sheet is variable depending on the basis weight of the paper sheet and the amount of the applied cationic compound, and is not limited to a specific value. Usually, when the paper sheet for the substrate sheet has a basis weight of 100 g/m2, preferably, the cationic compound-applied substrate paper sheet has a stxc3x6ckigt sizing degree of 1 to 100 seconds. If the sizing degree is too low, the application of the ink receiving layer may cause the substrate paper sheet to wrinkle. Also, if it is too high, the resultant ink jet recording sheet may exhibit an unsatisfactory ink absorption.\nOn the cationic compound-applied substrate sheet, an ink receiving layer comprising a pigment and a binder is coated. The ink receiving layer serves as an undercoat layer for the gloss layer.\nThe pigments usable for the ink receiving layer are selected from inorganic pigments, for example, kaolin, clay, calcined clay, amorphous silica (or non-crystalline silica), zinc oxide, aluminum oxide, aluminum hydroxide, calcium carbonate, satin white, aluminum silicate, alumina, colloidal silica, zeolite, sepiolite, smectite, synthetic smectite, magnesium silicate, magnesium carbonate, magnesium oxide and diatomaceous earth pigments, and organic pigments, for example, styrene polymer, urea resin, and benzoquanamine resin pigments, which are usually used for the production of coated paper sheets. Among the above-mentioned pigments, amorphous silica, aluminum oxide, zeolite and synthetic smectite pigments are specifically preferred because they can cause the resultant ink receiving layer to have a porous structure and exhibit an excellent ink absorption.\nIn another preferred embodiment of the present invention, the pigment contains amorphous silica and zeolite. When the zeolite is contained, the resultant recording sheet exhibits an enhanced image retention. Also, amorphous silica causes the resultant recording sheet to exhibit an enhanced color density of the images and an increased gloss of the gloss layer.\nThe zeolite is a general name of crystalline aluminosilicates and represented by the following general formula:\n(MI, MIIxc2xd)m(AlmSinO2(m+n))xc2x7XH2O\nwherein n and m respectively represent a numeral of 1 or more, and satisfy the relationship nxe2x89xa7m, MI represents an alkali metal ion, for example, Na+ and K+, MII represents an alkaline earth metal ion, for example, Ca2+, Mg2+ and Ba2+, and at least one of MI and MII must be present.\nThe zeolite is preferably selected from synthetic zeolite pigments which preferably have a particle size of 0.5 to 20 xcexcm.\nThe synthetic non-crystalline silica pigments are classified, in accordance with the production methods thereof, into wet process products and dry process products. In the wet process, a starting material consisting of sodium silicate is neutralized with an acid to cause fine silica particles to precipitate. The wet process is classified into precipitation method and gelation method. Also, in the dry process, the non-crystalline silica is produced by burning a starting material consisting of silicon tetrachloride with hydrogen and oxygen. The dry process non-crystalline silica pigment is called as white carbon, and known silicic anhydride or hydrated silicic acid. There is no limitation to the average particle size of the synthetic non-crystalline silica particles. Usually, the average particle size is preferably about 0.5 to about 20 xcexcm. For the purpose of enhancing absorption of the ink, preferably the non-crystalline silica particles have a specific surface area and an oil absorption as high as possible. More preferably, the specific surface area is 100 to 500 m2/g and the oil absorption is 100 to 400 ml/100 g.\nThe zeolite pigment is preferably used in a content of 3 to 60% by weight based on the total weight of the pigment.\nThe binder usable for the ink receiving layer preferably contains at least one member selected from, for example, proteins, for example, casein, soybean protein and synthetic proteins, starch compounds, for example, starch and oxidized starches, polyvinyl alcohols, cellulose compounds, for example, carboxymethyl-cellulose and methylcellulose, conjugated diene polymer latices, for example, styrene-butadiene copolymer and methyl methacrylate-butadiene copolymer latices, acrylic polymer latices, vinyl polymer latices, for example, ethylene-vinyl acetate latices, which are usually used for coating paper sheets. These compounds may be used alone or in mixtures of two or more thereof.\nThe binder for the ink receiving layer is preferably used in an amount of 1 to 40 parts by weight, more preferably 2 to 20 parts by weight, per 100 parts by weight of the pigment. If the amount of the binder is less than 1 part by weight, the resultant ink receiving layer may have an unsatisfactory mechanical strength and may be easily damaged or powdered. Also, if the binder is used in too high an amount, the resultant ink receiving layer may exhibit an insufficient ink absorption and may not exhibit a desired ink jet recording performance.\nThe ink receiving layer optionally contains at least one cationic compound selected from those the same as those usable for the substrate sheet. The cationic compound contained in the ink receiving layer contributes to enhancing water resistance of the ink images received in the ink receiving layer. There is no limitation to the content of the ink receiving layer. Usually, the cationic compound is contained preferably in an amount of 1 to 30 parts by weight, more preferably 5 to 20 parts by weight, per 100 parts by weight of the pigment, in the ink receiving layer.\nFurther, the ink receiving layer optionally contains an additive comprising at least one member selected from, for example, dispersing agents, thickening agents, anti-foaming agents, coloring agents, anti-static agents and preservatives. The ink receiving layer is formed by preparing a coating liquid for the ink receiving layer comprising the pigment, the binder and the optional cationic compound and/or other additive, in a total solid content of 5 to 65% by weight, coating a surface of the cationic compound-applied substrate sheet having a basis weight of 20 to 400 g/m2 with the coating liquid by using a conventional coating means, for example, blade coater, air knife coater, roll coater, brush coater, champlex coater, bar coater, gravure coater, die coater and curtain coater, and drying the coating liquid layer to form an ink receiving layer in a dry solid weight of 1 to 50 g/m2, more preferably 2 to 20 g/m2. Thereafter, optionally, the ink receiving layer is subjected to a smoothing treatment, for example, a super calender treatment and/or a brushing treatment.\nThe ink receiving layer is coated with a gloss layer. The gloss layer comprises a polymer resin and optionally a pigment. The polymer resin preferably consists of a polymerization product of at least one monomer having an ethylenically unsaturated bond. The gloss layer must allow the ink to easily pass therethrough and penetrate into the ink receiving layer. Namely, the gloss layer must be as porous and ink-permeable as possible unless the gloss of the gloss layer becomes unsatisfactory. For this purpose, it is preferable that the drying conditions for the coating liquid layer for the gloss layer are controlled so that the polymer resin incompletely forms a coating film, while creating a desired gloss. Also, when a pigment is contained, the pigment is preferably selected from those having a small particle size, for example, of about 300 nm or less preferably 200 nm or less, so as to not reduce the transparency and the gloss of the gloss layer. Such pigment can include colloidal silica, alumina sol and silica sol. The content of the pigment is adjusted to such an extent that the presence of the pigment does not cause a great reduction in the transparency of the resultant gloss layer. Usually, the pigment is contained preferably in a content of 200 parts by weight or less, more preferably 50 to 150 parts by weight, based on the 100 parts by weight of the polymer resin. If the content of the pigment is too much, the resultant gloss layer may exhibit an insufficient gloss.\nThe polymer resins for the gloss layer include, for example, polyvinyl alcohol resins, aqueous urethane resins and polymerization products of ethylenically unsaturated monomers.\nThe aqueous urethane resins can be produced by reacting polyisocyanate compounds, for example, diisocyanate, triisocyanate and tetraisocyanate compounds with polyol compounds.\nThe polymer resin for the gloss layer is preferably a polymerization product of at least one ethylenically unsaturated monomer. The ethylenicaily unsaturated monomers are preferably selected from acrylic esters having a hydrocarbon group with 1 to 18 carbon atoms, for example, methyl acrylate, ethyl acrylate, butyl acrylate, 2-ethylhexyl acrylate, lauryl acrylate, 2-hydroxyethyl acrylate and glycidyl acrylate; methacrylic esters having a hydrocarbon group with 1 to 18 carbon atoms, for example, methyl methacrylate, ethyl methacrylate, 2-hydroxyethyl methacrylate, 2-hydroxypropyl methacrylate, and glycidyl methacrylate; ethylenically unsaturated aromatic monomers, for example, styrene, xcex1-methylstyrene and vinyl toluene; aliphatic vinyl monomers, for example, acrylonitrile, vinyl chloride, vinylidene chloride, vinyl acetate, vinyl propionate, acrylamide and N-methylolacrylamide, and ethylenically unsaturated hydrocarbon monomers, for example, ethylene and butadiene.\nThe polymer resin for the gloss layer may comprise one or more selected from homopolymers of the above-mentioned monomers and copolymers of two or more of the above-mentioned monomers, and modification products of the homopolymers and copolymers. The modified polymers include, for example, carboxylated polymers and carboxylated polymers further modified to be alkali-reactive. Also, the modified polymer includes complexes produced by polymerizing the ethylenically unsaturated monomers in the presence of a colloidal silica so as to introduce Sixe2x80x94Oxe2x80x94R bonds (wherein R represents a group derived from the monomers) into the resultant polymers, or other complexes produced by introducing a functional group, for example, SIOH group, which is reactive with the colloidal silica, into the above-mentioned polymers or copolymers, and reacting the functional group-introduced polymers or copolymers with the colloidal silica.\nThe gloss layer is formed by coating a coating liquid containing, as a principal component, the polymer resin on the ink receiving layer, and drying the coating liquid layer.\nThe gloss layer may be formed by coating a gloss layer-forming coating liquid on a smooth surface of a plastic film, solidifying the coating liquid on the smooth plastic film surface, and transferring the resulting solid layer from the smooth plastic film surface to the ink receiving layer surface, to obtain a gloss surface of the gloss layer. The smooth surface of the plastic film preferably has a surface roughness Ra of 0.5 xcexcm or less, more preferably 0.05 xcexcm or less, determined in accordance with Japanese Industrial Standard (JIS) B0601.\nPreferably, the gloss layer is formed by coating the coating liquid containing the polymer resin on the ink receiving layer, pressing the coating liquid layer which is in wet condition, at a temperature of, for example, 50xc2x0 C. to 150xc2x0 C., preferably 70xc2x0 C. to 120xc2x0 C. against a mirror-finished drum surface, and drying the pressed coating liquid on the mirror-finished drum surface. Alternatively, the gloss layer is formed by coating the polymer resin-containing coating liquid on the ink-receiving layer, drying the coating liquid layer, rewetting the dried layer, and then pressing the rewetted layer against the mirror-finished drum surface and drying the pressed layer under the same conditions as mentioned above. These methods are referred to as a cast-coating method, and contribute to obtaining a gloss layer having excellent gloss, and ink-absorption. Among these methods, preferably the contact of the coating liquid layer with the mirror-finished casting drum surface is carried out while the coating liquid layer is still in a wet condition, namely by a wet cast-coating method.\nThe polymer resin for the gloss layer preferably has a glass transition temperature of 40xc2x0 C. or more, more preferably 50 to 100xc2x0 C. If the glass transition temperature is lower than 40xc2x0 C., when the coating liquid layer containing polymer resin is dried, the film formation of the polymer resin proceeds excessively, the porosity of the resultant gloss layer may be reduced too much, and thus the ink absorption of the gloss layer may become unsatisfactory.\nAlso, if the coating liquid layer for the gloss layer is dried at too high a drying temperature, the porosity of the resultant gloss layer may become too low and the ink absorption of the gloss layer may be unsatisfactory. Further, if the drying temperature is too low, the resultant gloss layer may exhibit an insufficient degree of gloss and the process efficiency in the gloss layer formation may be unsatisfactory. Accordingly, the drying procedure of the gloss layer preferably carried out at a temperature of 50 to 150xc2x0 C., more preferably 70 to 100xc2x0 C.\nThe gloss layer optionally contains an additive comprising at least one member selected from, for example, pigments, anti-foaming agents, coloring agents, anti-static agents, preservatives, dispersing agents and thickening agents which are usually used for the coated paper sheets for usual printing and ink jet recording, to adjust the brightness of the gloss layer and the viscosity and fluidity of the coating liquid, as desired.\nAs mentioned above, the coating liquid is coated on the ink receiving layer by using a conventional coating device, for example, a blade coater, air knife coater, roll coater, brush coater, champlex coater, bar coater or gravure coater, die coater or curtain coater. Preferably, the coating liquid layer is pressed, in wet condition, against a mirror-finished surface of a casting drum at a temperature of 50 to 150xc2x0 C., and dried so as to transfer the mirror surface to the surface of the gloss layer. The gloss layer is preferably formed in a dry solid amount of 0.2 to 30 g/m2, more preferably 1 to 10 g/m2. If the amount is too small, the resultant gloss layer may exhibit an unsatisfactory gloss. Also, if the amount is too great, the resultant gloss layer may cause an unsatisfactory ink drying rate and a reduced color density of printed ink images. Optionally, the gloss layer may be smoothed by a super calender treatment.\nThe ink jet recording sheet of the present invention exhibits excellent gloss and superior ink jet recording performance, due to the fact that the cationic compound is applied to the substrate sheet. The mechanism of enhancing the gloss and ink jet recording performance by the cationic compound applied to the substrate sheet is not completely clear. It is assumed, however, that when a coating liquid for the ink receiving layer is coated on a surface of a substrate sheet impregnated or coated with a cationic compound, the coagulation of the coating liquid is promoted by the cationic compound so that an excessive penetration of the coating liquid into the substrate sheet is prevented, and thus an ink receiving layer with an appropriate porosity is formed. Then, when a coating liquid for the gloss layer is coated on the ink receiving layer, penetration of the coating liquid into the ink receiving layer and then into the substrate sheet is restricted for the same reasons as mentioned above, and therefore a gloss layer having a high and uniform gloss is formed. The inventors of the present invention have confirmed that when the cationic compound was omitted from the ink jet recording sheet of the present invention, or replaced by a non-cationic substance, for example, binder, the specific effects of the present invention could not be obtained.\nIn the ink jet recording sheet of the present invention, the above-mentioned specific effects of the present invention can be further enhanced by adding a cationic compound to the ink receiving layer.\nThe ink jet recording sheet of the present invention has a gloss layer with excellent gloss and an ink receiving layer which is porous and thus exhibits a high ink absorption, and a high resistance to surface-roughening due to swelling of portions of the ink receiving layer by the received ink. Therefore, as a whole, the ink jet recording sheet can maintain the high gloss of the gloss layer surface even after printing.\nThe ink jet recording sheet of the present invention enables the recorded ink images to exhibit a high and uniform color density and a high water resistance."}
{"text": "If a printing job is generated when an image forming apparatus is in a sleep state, the image forming apparatus returns from the sleep state and performs a printing preparation operation (hereinafter referred to as pre-run). A period from the start of the pre-run until printing is actually executed mainly depends on time for warming a fixing unit (hereinafter, heat accumulation time).\nIn the past, when the image forming apparatus returns from the sleep state, irrespective of what kind of a printing job is performed next, the image forming apparatus always warms the fixing unit until a fixed time elapses."}
{"text": "1. Field of Invention\nThe present invention relates to a fan system, and more particularly to a fan system capable of storing energy by itself for re-usage.\n2. Related Art\nAccompanying to the development of technology, the functionality of electronic system is upgraded and the integration of electronic devices is also raised. Electronic system is generally disposed at least one fan for heat dissipation to keep electronic system being normally operated.\nIn order to enhance the effect of heat dissipation, a plurality of fans are connected with each other in series to raise the pressure and volume of airflow. As shown in FIG. 1, a conventional fan system 1 includes a first fan 11 and a second fan 12. An external voltage source P is connected with the first fan 11 and the second fan 12 when the first fan 11 and the second fan 12 are connected in series. However, the second fan 12 is rotated by the first fan 11 and needs not to be driven by the voltage source P when the first fan 11 is driven to rotate by the voltage source P. In addition, the second fan 12 is driven by the voltage source P when the first fan 11 is out of order. The first fan 11 is thus driven to rotate by the second fan 12. Under such construction, the first fan 11 or the second fan 12 is driven by the other fan, and the energy generated by the driven fan is wasted and not recycled for re-usage. The heat dissipation of the electronic system will be affected if the voltage source P is out of order.\nIt is thus imperative to provide a fan system capable of storing energy by itself for re-usage."}
{"text": "With economic and social globalization, an increasing quantity of enterprises and organizations establish branches nationwide and even worldwide, and an increasing quantity of employees on business trips access headquarters' service systems from fixed locations. In addition, as services continuously expand and application becomes increasingly complex, an information technology (IT) infrastructure of an organization is developing toward centralized management, and setting up data centers and cloud services has been a trend of IT construction. An enterprise service gateway is an entrance to a data center and cloud service center. A throughput capability of the enterprise service gateway directly affects performance of the entire data center and cloud service center.\nOPENFLOW, a new network switching model supporting network innovation researches, was put forward by Stanford University in the United States. The OPENFLOW protocol is gradually improved and successfully applied to an actual network, where application of the OPENFLOW protocol to a data center network is particularly prominent. As shown in FIG. 1, a topology structure of an OPENFLOW network includes an OPENFLOW controller 11, an OPENFLOW switch 12, a terminal 13, and links connecting the OPENFLOW controller 11, OPENFLOW switch 12, and terminal 13. The OPENFLOW controller 11 stores a topology structure of an entire network, generates a flow table for a data flow that needs to be forwarded, and delivers the flow table to a corresponding switch. The OPENFLOW switch 12 stores flow table information delivered by the OPENFLOW controller 11 and performs data forwarding between the OPENFLOW switches 12 and between the OPENFLOW switch 12 and the terminal 13.\nA topology structure of the data center network is generally a tree structure. The enterprise service gateway, as the entrance (a root) to the data center, is responsible for network forwarding and further needs to perform particular processing on data, for example in terms of network security, to perform “encryption/decryption, Deep Packet Inspection (DPI), and the like” on the data, and to perform data compression and decompression, and the like. The processing includes the following steps:\n1. The enterprise service gateway receives a client request and performs data processing such as security detection (Internet Protocol Security (IPSEC), and DPI) on data.\n2. After the data processing is completed, a to-be-processed packet is forwarded to a back-end server according to a service type of the request.\n3. The server performs processing, and forwards a result of the processing to the enterprise service gateway.\n4. The enterprise service gateway receives service data, performs particular processing on the service data, and forwards an encapsulated packet to a client.\nAn enterprise service gateway serves as an entrance to an entire data center. When a to-be-processed data volume exceeds a capability of the enterprise service gateway, the enterprise service gateway becomes a bottleneck of an entire system, whereas a large quantity of computing resources on a back-end server are idle. Therefore, a method for resolving the bottleneck problem needs to be found."}
{"text": "1. Field\nMethods and apparatuses consistent with exemplary embodiments relate to an image display apparatus and an image display method, and more particularly, to an image display apparatus and an image display method capable of substantially saving power consumption.\n2. Description of the Related Art\nCathode ray tubes (CRTs) have been widely used in display apparatuses such as televisions (TVs) and computer monitors. However, CRTs use a method of emitting cathode electrons from an electron gun to glass wool covered with a fluorescent material to generate light, and thus have large sizes, which is disadvantageous. Thus, a flat panel display (FPD) technology emerged, which drastically reduces the thickness of a display apparatus.\nExamples of flat panel displays include a liquid crystal display (LCD) that uses a liquid crystal, a plasma display panel (PDP) that emits light by using a principle of gas discharge such as that employed in neon signs, an organic light emitting diode (OLED) panel that emits light by applying a current to a fluorescent organic compound, a field emission display (FED), an electro luminescence display (ELD), and so forth.\nVarious technologies are proposed to reduce power consumption in flat panel displays. However, as the size of display panels increases and the resolution becomes higher, the power saving technologies proposed in the related art have limits to substantially reducing power consumption. Therefore, there is a need for a method that may substantially reduce power consumption in a display apparatus."}
{"text": "When transparent glass panes first became available for widespread use by consumers, large panes of such glass were prohibitively expensive and difficult to form. As a result, early windows and doors were constructed of a plurality of small, inexpensive, easily manufactured glass panes. Early windows and doors were characterized by a grid-like cross-hatched appearance as a result. The windows and doors of many old houses are a testament to this early technology.\nGlass manufacturing technology rapidly advanced, however, making possible the introduction of the \"picture window\", a name sometimes applied to large panes of glass. While large uninterrupted panes of glass gained substantial popularity as early as the 1940's, doors and windows having multiple glass panes are considered by many to be more aesthetically pleasing. The irony of technology, however, is that the earlier design of windows and doors consisting of multiple panes of glass separated by dividing elements is now much more expensive to produce than a single pane window.\nIn order to meet the demand for doors and windows having the appearance of older style, multiple-pane assemblies, a number of techniques have been adopted to give an aesthetically pleasing appearance to a single large pane of glass.\nIt is commonly known to produce a unitary latticework of wood, metal or plastic, and apply this latticework to one side of a large contiguous sheet of glass. From a distance, the ordinary observer sees the assembly as comprised of multiple panes of glass constructed and assembled by the latticework. In fact, if two identical such latticework appliques are placed on opposite sides of a single pane of glass, it is difficult, without close inspection, to perceive that the assembly is, in fact, a decorative simulation of multiple panes of glass.\nIn recent years, with the escalating cost of energy, more and more buildings have incorporated insulated glass windows and doors in place of earlier single pane windows and doors, which are notoriously poor insulators. Modern insulated glass consists of multiple glass sheets, separated by spacers at the perimeter of the sheets, and then sealed together with a semi-flexible rubber-like compound. Such a sandwich of glass and air provides an enclosed insulating air space, which results in desirable thermal insulating properties.\nDouble-paned assemblies, however, are not easily disguised by latticework appliques to simulate the appearance of multi-panel windows and doors. Because there is a space of approximately 3/8\" between the glass panels, the application of a latticework applique is not aesthetically pleasing. A more suitable visual effect is obtained by inserting a latticework assembly between the two glass panes. Windows and doors containing such internal latticeworks have become very popular in modern commercial and residential buildings.\nThe construction of an insulated window using such internal grid work is, however, a process which is relatively work-intensive. Because there are an infinite number of sizes and shapes of windows and doors in use today, the manufacture of grid work for glass panes dictates the availability of an equally infinite variety of gridwork assemblies. This problem is further compounded by the popularity of two distinct gridwork styles: those which are primarily rectangular in appearance, i.e., having essentially vertical and horizontal members; and those which have a diamond appearance, where the intersecting members of the grid are inclined at angles which are not strictly horizontal or vertical.\nThe manufacture of each individual gridwork element, therefore, requires several steps. First, the type of gridwork pattern to be incorporated in the window or door must be determined. Further, the dimensions of the finished insulated glass pane must be determined. The final determination is the spacing of the various intersection gridwork elements. With these parameters, it is possible to cut the individual intersecting elements of the grid to their appropriate lengths. Finally, the intersecting elements must be notched along their length, so that the completely assembled gridwork will have a uniform thickness, measured from front to back, to fit within the appropriate insulated glass insulation space.\nThe material which is most commonly used for the construction of this gridwork is an extruded aluminum tube of rectangular cross-section, which is painted or clad in an appropriate color. Originally, the manufacture of the gridwork required the determination steps above-described, followed by manual cutting of the extruded aluminum element from an elongated piece of stock material, known as a muntin bar. Muntin bar is typically supplied in lengths of 144\", providing enough stock material for several individual grid elements.\nIn the existing technology, the manufacture of the gridwork is done essentially by hand, with manual computation of the lengths of grid elements required to be cut from the available lengths of muntin bar. Likewise, manual computations of the locations for notching the grid elements are performed. It is important to note that when manufacturing diagonally-oriented of grid elements, it is frequently necessary to insure that the ends of the elements are cut at the appropriate angle, so that the ends will fit flush with the interior surface of the spacer which separates the panes of the insulated glass. This additional requirement further complicates the manufacture of diagonally oriented grid elements.\nRecently, a number of computer programs have become available which assist the workmen in the cutting and notching of grid elements, by providing information regarding the angles of the ends of the grid elements, as well as the precise location for the notches where the grid elements will fit together. Still, the cutting and notching operations take place manually, requiring the workmen to move the raw stock muntin bar into position against either a cutting saw or router, and positioning the angle of the cutting saw or router manually, based on the computations made either manually, or by the computer program. Modern systems, such as the muntin notcher and muntin angle and trim saw produced by McKeegan Equipment & Supply Company of Plymouth, Michigan, provide minimal digital readout settings, and manual stops for manual positioning of the muntin bar at the cutting or notching stage of a machine. These processes, even though improved, are still work-intensive, time-consuming and expensive. The present invention overcomes each of these limitations."}
{"text": "Licensed wireless systems provide mobile wireless communications to individuals using wireless transceivers. Licensed wireless systems refer to public cellular telephone systems and/or Personal Communication Services (PCS) telephone systems. Wireless transceivers include cellular telephones, PCS telephones, wireless-enabled personal digital assistants, wireless modems, and the like.\nLicensed wireless systems utilize wireless signal frequencies that are licensed from governments. Large fees are paid for access to these frequencies. Expensive base station equipment is used to support communications on licensed frequencies. Base stations are typically installed approximately a mile apart from one another. As a result, the quality of service (voice quality and speed of data transfer) in wireless systems is considerably inferior to the quality of service afforded by landline (wired) connections. Thus, the user of a licensed wireless system pays relatively high fees for relatively low quality service.\nLandline (wired) connections are extensively deployed and generally perform at a lower cost with higher quality voice and higher speed data services. The problem with landline connections is that they constrain the mobility of a user. Traditionally, a physical connection to the landline was required. Currently, unlicensed wireless communication systems are deployed to increase the mobility of an individual using a landline. The mobility range associated with such systems is typically on the order of 100 meters. A common unlicensed wireless communication system includes a base station with a physical connection to a landline. The base station has a RF transceiver to facilitate communication with a wireless handset that is operative within a modest distance of the base station. Thus, this option provides higher quality services at a lower cost, but the services only extend a modest distance from the base station.\nThus, there are significant shortcomings associated with current landline systems and licensed wireless systems. For this reason, individuals commonly have one telephone number for landline communications and one telephone number for licensed wireless communications. This leads to additional expense and inconvenience for an individual. It would be highly desirable if an individual could utilize a single telephone number for both landline communications and licensed wireless communications. Ideally, such a system would allow an individual, through seamless handoffs between the two systems, to exploit the benefits of each system."}
{"text": "Catalytic perfluoroalkylation of aromatic compounds is known in the art. The process, as taught in the art, provides a method whereby an unfunctionalized aromatic group is substituted with a perfluoroalkyl group with the concurrent production of hydrogen iodide. A process for the catalytic addition of perfluoroalkyl iodide to benzene using a ruthenium catalyst on carbon and employing solid potassium carbonate as a hydrogen iodide acceptor is disclosed in European Patent Application No. 0 114 359 of Werner. Good yields of the perfluoroalkyl benzene are reported but only when large amounts of catalyst are used."}
{"text": "Cancer is a disease resulting from an abnormal growth of tissue. Certain cancers have the potential to invade into local tissues and also metastasize to distant organs. This disease can develop in a wide variety of different organs, tissues and cell types. Therefore, the term “cancer” refers to a collection of over a thousand different diseases.\nAnaplastic lymphoma kinase (ALK), a member of the insulin receptor superfamily of receptor tyrosine kinases, has been implicated in oncogenesis in hematopoietic and non-hematopoietic tumors. The aberrant expression of full-length ALK receptor proteins has been reported in neuroblastomas and glioblastomas; and ALK fusion proteins have occurred in anaplastic large cell lymphoma. The study of ALK fusion proteins has also raised the possibility of new therapeutic treatments for patients with ALK-positive malignancies. (Pulford et al., Cell. Mol. Life Sci. 61:2939-2953 (2004)).\nInsulin-like growth factor (IGF-1) signaling is highly implicated in cancer, with the IGF-1 receptor (IGF-1R) as the predominating factor. IGR-1R is important for tumor transformation and survival of malignant cells, but is only partially involved in normal cell growth. Targeting of IGF-1R has been suggested to be a promising option for cancer therapy. (Larsson et al., Br. J. Cancer 92:2097-2101 (2005)).\nThe c-ros oncogene 1 (ROS1, also known as ROS), a member of the tyrosine kinase insulin receptor gene family, is highly expressed in a variety of tumor cell lines.\nDespite advancements in the art, there remains a need for cancer treatments and anti-cancer compounds."}
{"text": "1. Field of the Invention\nThe present invention relates to a solar collector for generating electricity and heat, and more particularly to a hybrid solar collector for converting solar energy of short wavelength into electric energy and converting solar energy of long wavelength into heat. The solar rays are separated into short wavelength and long wavelength. Separated short wavelength ray is focused on a solar cell for converting the energy of the short wavelength into electric energy, and separated long wavelength ray is focused on a pipe for heating a liquid inside the pipe by the energy of the long wavelength ray.\n2. Description of Related Art\nAs a solar energy collector for generating electric energy, a solar panel is well known. The solar panel is constituted of solar cells arranged in planar form. Further, a large-scale solar collector is also known. Such a solar collector is constituted of a large number of reflector mirrors distributed on the ground for converging solar rays at a converging point, and a solar cell arranged at the converging point.\nOn the other hand, a solar collector for generating heat is well known. Such a solar collector is constituted of reflector mirrors for converging solar rays, and a heat absorbing tube arranged at a converging point.\nMoreover, another solar collector for generating electricity and heat has been proposed. The hybrid solar collector utilizes a solar cell arranged at a solar rays converging point and a heat absorbing tube located on a rear surface of the solar cell."}
{"text": "The various embodiments and aspects described herein relate to an accessory for a truck bed to elongate a length of the truck bed.\nTrucks have beds that come in a variety of sizes. By way of example and not limitation, the truck beds may be a short bed or long bed. The short bed may serve the purposes of its owner a majority of the time but at certain times, the owner may require a slightly longer truck bed to haul larger or longer objects. In this regard, the owner may let down his or her tailgate and allow the objects to hang out of the back side of the truck bed when hauling the long object. In doing so, the object becomes a safety hazard in that other vehicles and pedestrians may not see the elongated object hanging out of the truck bed and may inadvertently hit the object causing injury to the person and/or damage to the object.\nAdditionally, the objects configuration may not be conducive to easily tying the object down in the truck bed. The tailgate is typically lifted upward in order to trap the object within the truck bed so that even if the object slides in the truck bed, the object would still remain within the truck bed.\nAccordingly, there is a need in the art for an accessory for a truck bed."}
{"text": "In recent years, logically dividing one physical computer into plural computers using a virtual technology, these computers respectively run a separate Operating Systems (OS), for example, a guest OS. These logical computers may be referred to as a virtual machine (VM).\nIn addition, according to the virtual technology, consolidating or integrating the OS running on plural servers into one server, it may be possible to realize power saving by efficient use of resources. Further, a migration technology for dynamically shifting the running virtual machine to other server has been also known. Depending on this technology, a power-saving operating method for cutting power of redundant servers by efficiently consolidating the virtual machines in accordance with workload has been also known.\nOn the other hand, the virtual technology originally for a large computer such as a server started to be applied to client's devices such as a notebook PC and a cellular phone because of downsizing and significant increase in performance of the computer. Due to application of the virtual technology to these client's devices, it is possible to expect improvement of a convenience obtained by using plural guest OSs with different user interfaces and achievement of security depending on abstraction of hardware and management the OS from a lower layer. Further, the notebook PC and the cellular phone are directly used by a user, so that power-saving operation used by a server for consolidating or migrating one device to other is hardly applied to these devices.\nAn information processor such as a notebook PC and a cellular phone to be mainly driven by a battery highly requires reduction of power consumption, and conventionally, a state of a power source of a device has been carefully managed by the OS and applications. For example, such an information processor realizes a function such as shifting the device into a lower power consumption mode, for example, clock-down and power supply stop. However, it is not possible to apply these methods of managing a power source state to a virtual machine system because there is a possibility that plural guest OSs are carried out at the same time and the same device is shared by other guest OS, so that this makes one guest OS impossible to freely stop the operation state.\nTherefore, conventionally, a method of arranging a state transition command notified by a guest OS (for example, Advanced Configuration and Power Interface (ACPI) command) by means of a management OS and managing device resources shared by this management OS has been known, for example. Thereby, controlling shift to a real power saving mode of the device, it is possible to reduce power consumption in the virtual machine system in which plural guest OSs run.\n[Patent Document 1] Japanese Patent Application Publication No. 2005-115653\n[Patent Document 2] Japanese Patent Application Publication No. 2006-113767\nHowever, as described above, according to the method of arranging a state transition command from each guest OS running in the virtual machine by means of the management OS, in order to appropriately realize power saving by the entire system, it is required that all guest OSs appropriately issue the state transition commands at anytime to positively carry out power saving. In other words, this involves a problem that, if only one guest OS does not issue the state transition command, shift to the power saving mode cannot be realized.\nOn operation of the system, a legacy OS unequipped with a function to command a state transition command may be run as a guest OS in order to maintain a compatibility. In such a case, for example, it is difficult to modify the legacy guest OS so as to be able to issue an appropriate state transition command. In addition, not only for the legacy guest OS but also for all guest OSs, it may be difficult to have equal power saving mechanisms due to a development resource problem.\nOn the other hand, under the condition that the guest OS does not run actively, the device maybe shifted to the power saving mode not depending on the transition command from the guest OS but in the judgment of the management OS and the virtual machine monitor. This may have a problem that it takes time since the device should be recovered from the power saving mode when the guest OS tries to use the device."}
{"text": "(i) Technical Field\nThe present disclosure relates to a semiconductor memory device and to a method of fabricating the same, and more particularly, to a semiconductor memory device including a cylinder type storage node and to a method of fabricating the semiconductor memory device.\n(ii) Description of the Related Art\nA memory cell of a direct random access memory (DRAM) among semiconductor memory devices may include a field effect transistor as a switching device for controlling read/write operations, and a capacitor for storing information. As the scale of the DRAM has been reduced, the area occupied by the capacitor of the memory cell has been continuously decreased. To ensure a sufficient effective capacity despite the reduction of cell area, a concave type or a cylinder type storage node having a three-dimensional capacitor shape has been suggested. The cylinder type storage node is widely used in DRAM in association with a capacitor over bit line (COB) structure, in which the capacitor is formed on a bit line.\nThe cylinder type storage node may be required to have a height of about 2 μm or higher to obtain an effective capacity of about 25 fF or greater per unit memory cell in the case of one cylindrical storage node (OCS) structure with a design rule of about 70 nm. As the cylinder type storage node has a very large aspect ratio, the cylinder type storage node may be structurally unstable. The unstable structure of the cylinder type storage node may cause defects such as, for example, slanting or lift off of the cylinder type storage node, which may cause a twin bit error or a multi-bit error during the operation of the DRAM.\nIn general, to ensure the mechanical stability of the cylinder type storage node which has a large aspect ratio, an etch stop layer may be formed under a lower electrode of the cylinder type storage node for supporting the cylinder type storage node. However, as the aspect ratio of the cylinder type storage node gradually increases, there may be a limit in providing the cylinder type storage node with sufficient mechanical stability by using the etch stop layer. In addition, when an etching process is performed to form a contact hole which exposes a storage node contact plug in the etch stop layer, it may be difficult to control critical dimensions (CD) on a lower portion of the cylinder type storage node due to skew of an etched surface of the etch stop layer."}
{"text": "As is well known in the field of integrated circuit design, layout and fabrication, the manufacturing cost of a given integrated circuit is largely dependent upon the chip area required to implement desired functions. The chip area, in turn, is defined by the geometries and sizes of the active components such as gate electrodes in metal-oxide-semiconductor (MOS) technology, and diffused regions such as MOS source and drain regions and bipolar emitters and base regions. These geometries and sizes are often dependent upon the photolithographic resolution available for the particular manufacturing facility. The goal of photolithography in establishing the horizontal dimensions of the various devices and circuits is to create a pattern which meets design requirements as well as to correctly align the circuit pattern on the surface of the wafer. As line widths shrink smaller and smaller in submicron photolithography, the process to print lines and contact holes in photoresist becomes increasingly more difficult.\nWith circuit advancement to the ultra-large-scale integration (ULSI) levels, more and more layers are added to the surface of the wafer. These additional layers in turn create more steps on the wafer surface. The resolution of small image sizes in photolithography thus becomes more difficult over the additional steps due to the increased problem of depth of focus. Planarization techniques become increasingly more important to offset the effects of a varied topography.\nPlanarization techniques may be applied to both dielectric layers and conductors or semiconductors. Planarizing regions of a wafer may be done in degrees from smoothing a particular layer to reduce the severity of steep slopes to global planarization across the entire wafer regardless of the underlying topography. Of course, global planarization is ideal but is also more difficult to achieve particularly between widely isolated features. The variations in the thickness of the material to be etched, differences in adjacent materials being etched and the underlying topography all add to the degree of difficulty of achieving global planarization.\nInterconnect technology, creating connections between conducting regions, relies heavily on planarization techniques. In modern integrated circuits, the material of choice for upper-level conductive interconnection systems has been aluminum, including doped aluminum and aluminum alloys. Aluminum is an attractive material for integrated circuit metallization due to its high conductivity and low cost. The processing required to form aluminum metallization is also relatively easy, as it can readily be evaporated or sputtered onto the wafer. Aluminum is also able to form good ohmic contact to both p-type and n-type doped semiconductor material, such as silicon. In addition, aluminum is quite compatible with conventional semiconductor processes, such as that used to form bipolar and metal-oxide-semiconductor (MOS) devices, unlike other metals such as copper or gold which can diffuse into active regions and degrade device performance.\nCertain drawbacks do exist for aluminum-based metallization systems, however, particularly as geometries enter the submicron regime. A well-known limitation of aluminum is its poor step coverage, particularly for vertical or retrograde sidewalls of contact openings through insulating layers such as silicon dioxide, and especially for sputtered aluminum, due to the shadowing effect of steep contact walls. In addition, mechanical stress in the aluminum film can cause voids therein. Stress-induced voids and step coverage faults of sufficient size can each cause an open in a metal line or contact. Furthermore, since aluminum metallization is subject to electromigration, and since the rate of electromigration increases with current density through the film, necking or narrowing of an aluminum line due to such voids or poor steps locally increases the current density. As a result, the electromigration rate increases at narrowed locations of the film, greatly increasing the electromigration failure rate.\nA prior technique for addressing these limitations of aluminum metallization systems includes the use of refractory metal plugs, such as tungsten plugs, to fill contact openings in insulating layers. According to one example of this technique (for which many specific methods are well known in the art), after the opening of contacts through the insulating layer, a layer of tungsten is deposited by CVD over the wafer in such a manner as to conformally fill the contact opening, and is subsequently etched back to expose the surface of the insulating layer with the tungsten remaining in the contact opening. Alternatively, selective tungsten deposition has been used (the tungsten deposition on silicon but not on silicon dioxide) to fill contact openings. In either case, a subsequently deposited aluminum layer can readily make contact to the tungsten plug.\nWhile the tungsten plug technique has many advantages, including good step coverage in all contacts, compatibility with planarized processing, and tolerance of misalignment in the etch of overlying aluminum lines (since the aluminum can be etched selectively relative to the tungsten), the tungsten plug process adds complexity to the manufacturing flow. In addition, deposited tungsten is vulnerable to poor adhesion and high contact resistance, requiring the use of additional sputtered barrier films prior to the deposited tungsten, and the associated added process complexity therewith.\nAs circuit density and device performance requirements increase, more complex wiring or routing of interconnects between conductors is required. Vertical as well as horizontal interconnects are required as the number of metallization layers increase to meet the requirements of more complex circuits. Current generation devices may require up to 5 to 6 layers of metallization to meet the wiring needs of state-of-the-art devices such as logic devices and SRAMs. To make such devices manufacturable within the limitations of existing equipment, planarization at each stage of processing becomes even more critical.\nAccordingly, it is an object of the present invention to provide a method of forming planarized multilevel metallization regions to allow for complex routing of metal interconnections to achieve greater device density.\nIt is a further object of the present invention to form planarized regions containing metal interconnections and vertical metal pillars or vias between metal interconnections to support the multilevel metallization process.\nIt is yet a further object of the present invention to form planarized aluminum interconnections and vias in submicron geometries with adequate step coverage absent of stress-induced voids and electromigration problems.\nIt is a further object of the present invention to provide such a method which utilizes standard processing techniques.\nIt is yet a further object of the present invention to provide such a method which allows for closer spacing of interconnect lines enabling more devices to be formed in a smaller chip area.\nOther objects and advantages of the present invention will be apparent to those of ordinary skill in the art having reference to the following specification together with the drawings."}
{"text": "1. Field of the Invention\nThis invention in general relates to a method for forming semiconductor devices and more particularly, to a method for forming shallow trench isolations (STI) structure on a semiconductor substrate.\n2. Description of the Related Art\nIsolation structure on a semiconductor substrate is used to prevent carriers from migrating to the adjacent devices. Typically, the isolation structure is formed in a semiconductor substrate with densely packed devices, such as in DRAM or in field effect transistor (FET) devices, for mitigating leakage current induced along the edges of EFT devices.\nOne method for forming the isolation structure in the semiconductor substrate is local oxidation. However, there are some inherent problems in local oxidation technology, such as the stress of the isolation structure or bird's beak encroachment incurred thereby. Therefore, in high-density semiconductor devices, the use of shallow trench isolation structure is proposed.\nThe method for forming shallow trench isolation structures in semiconductor substrate comprises the application of anisotropic etching on a semiconductor substrate to form trenches therein and then filling the trenches with dielectric material. Since the field isolation effect provided by the shallow trench isolation structures is scaleable and the Bird's Beak Encroachment in field oxidation is also resolved, the technology of shallow trench isolation structure has become a trend in the technology of semiconductor.\nFIGS. 1A to 1B schematically illustrate in cross-sectional representation of the conventional method for forming a shallow trench isolation structure. Referring to FIG. 1A, a pad oxide layer 102 and a silicon nitride layer 104 are provided on a substrate 100. The silicon nitride layer 104 is patterned as a mask to etch a trench 106 in the substrate 100. Then, a liner oxide layer 108 is formed on the surface of substrate exposed in the trench 106 by dry oxidation. Next, a silicon oxide layer 110 is overlaid on the substrate 110 and fills the trench 106 therein. Now referring to FIG. 1B, the oxide layer 110 is removed by chemical mechanical polishing until the silicon nitride layer 104, which is used as a stop layer, is exposed. Next, the silicon nitride layer 104 is removed and then, the pad oxide layer 102 is removed by hydrofluoric acid. A trench isolation structure 112 is formed to completion in the substrate 100.\nIn the above-described method, the silicon oxide layer 110 covered on the substrate is formed by chemical deposition, while the pad oxide layer is formed by thermal oxidation. The compactness of the former is less than that of the latter. Therefore, while using hydrofluoric acid as an etchant to remove the pad oxide layer 102, the silicon oxide 110 in the trench 106 is etched at a rate higher than that for pad oxide layer 102. As a result, lateral etching occurs on the top surface of trench 106 to form a groove 116, which results in the less-smooth corner profile at the corner 114. Due to the lateral etching, the gate layer subsequently formed on the substrate will be thinner at the corner area 114 than in other places, so the kink effect will occur on the substrate.\nTherefore, a need exists to improve the method for forming shallow trench isolations (STI) structure on a semiconductor substrate to eliminate the above mentioned problems."}
{"text": "This invention relates generally to returnable beverage can containers.\nIn order to help with the problem of waste disposal, many states have made laws requiring retail stores to receive empty cans and bottles and to refund the deposit on each. The most common type of cans returned are 16 ounce cylindrical cans having substantially the same general diameters and lengths.\nContainers have been devised to store and hold the cans. A problem of these containers is that a number of the empty cans are accumulated and stored prior to transporting them to a retail store for the return of the deposit. Empty beverage cans generally have small top openings and are therefore difficult to clean well and so tend to attract vermin, primarily various types of insects.\nA number of U.S. patents that have been issued that describe various types of empty beverage can containers are as follows:\n(1) U.S. Pat. No. 4,214,660 issued to Hunt on July 29, 1980. Hunt describes an empty can container that is analogous to the packages for holding full cans on sale such as the six or twelve pack container. Inspection holes are positioned over each can.\n(2) U.S. Pat. No. 4,290,525 issued to Sisson on Sept. 22, 1981. Sisson describes a transparent carrier for empty beverage cans that is constructed so as to show some of the cans for inspection and allows quick counting of all of them. This invention seems to be primarily directed to the container having a base stiffener panel and pleated end panels to allow the container to be folded around the base panel for storage of the container.\n(3) U.S. Pat. No. 4,299,324 issued to Dickens on on Nov. 10, 1981. Dickens describes a collapsible compartmented container made of sheet plastic for returnable beverage cans.\n(4) U.S. Pat. No. 4,417,657 issued to Thibodeau on Nov. 29, 1983. Thibodeau describes an open top container that has chambers for holding the can in vertical stacks.\n(5) U.S. Pat. No. 4,542,826 issued to Adams on Sept. 24, 1985. Adams describes a flexible, transparent bag having vertical compartments for receiving beverage cans.\n(6) U.S. Pat. No. 4,574,978 issued to Hodges on Mar. 11, 1986. Hodges describes a can container that includes a transparent, flexible pouch and a rigid housing in which the pouch is slidably received. Hodges further describes a removable cover for covering the container during transportation and also a flap for covering the pouch during storage.\nNone of the described containers is capable of keeping out vermin so as to prevent them from entering the empty cans during the storage phase of the empty cans."}
{"text": "Examples of known fuel injection systems use an injector to dispense a quantity of fuel that is to be combusted in an internal combustion engine. The quantity of fuel that is dispensed is varied in accordance with a number of engine parameters such as engine speed, engine load, engine emissions, etc.\nKnown electronic fuel injection systems monitor at least one of the engine parameters and electrically operate the injector to dispense the fuel. It is believed that examples of known injectors use electromagnetic coils, piezoelectric elements, or magnetostrictive materials to actuate a valve.\nA known fuel injector utilizes a plethora of internal components such as a metallic inlet tube connected to a valve body via a non-magnetic shell with a pole piece interposed therebetween. The inlet tube, valve body, non-magnetic shell and pole piece are generally affixed to each other after a closure assembly and a metering assembly are disposed in the valve body. A solenoid coil is inserted over the assembled components and the entire assembly is molded into the fuel injector.\nIt is believed that one known fuel injector utilizes a plastic body molded over a solenoid coil to provide a plastic inlet fuel passage with a metallic valve body being coupled to the solenoid coil.\nIt is believed that another known fuel injector utilizes two separate subassemblies to form the fuel injector. The first subassembly can include a complete coil assembly and electrical connector molded into an outer casing to provide a power group. The second subassembly can include an inlet tube, pole piece, non-magnetic shell valve body, closure assembly and metering assembly affixed together to form a stand alone fuel group. The two sub-assemblies are formed separately and coupled together to provide an operable fuel injector.\nWhile the known fuel injectors are suited to the task of metering fuel, it is believed that the known fuel injectors may have certain assembly or component drawbacks that require extensive manufacturing process to be undertaken to ensure that the injector are suitable for commercial applications. They can include, for example, the necessity for multiple seal points between components to provide leak integrity in the injector and a large number of manufacturing steps that are undertaken. These seals can be effectuated by elastomeric seals, such as, O-rings, or multiple hermetic welds to ensure structural and leak integrity of the known fuel injectors. Others include the potential manufacturing difficulties associated with thermal distortion in welding multiple metallic components at close proximity to each other or the need for a metal valve body with internal resilient seals for leak integrity. Yet another drawback can include the utilization of lift setting components that must be inserted into the valve body of the fuel injector. Thus, it would be advantageous to reduce or even eliminate some of these drawbacks."}
{"text": "The invention relates to a device for limiting the extent of movement of a moveable member and is particularly suitable in the field of focusing mechanisms for cameras, and in particular is directed to a device for cutting off a motor which drives a focusing member when the focusing member arrives at a position determined by a mechanical stop member.\nIn many cameras, such as television cameras, movie cameras or the like it is common to provide a motor driven objective focusing lens. Such may also appear in still cameras having automatic focusing. The lens is mounted on a moveable member, such as a rack of a rack and pinion mechanism, and, a motor is energized to move the moveable member, such as, by rotating the pinion. It is also common in such devices to provide a mechanical stop member for insuring that the moveable member is accurately stopped at a precise position for infinity object focusing. However, there are disadvantages attendent to the mechanical stop member.\nWhen the focusing lens moveable member impinges on the stop member, an excessive force is exerted on the electric motor possibly harming or causing failure of the motor. If a battery is used as the energy source, there is the additional disadvantage of unnecessary battery discharge at this time.\nOne method for overcoming the latter problems is to provide a clutch mechanism between the motor and the moveable member. This has been found to be unsatisfactory because even after interruption of the movement of the moveable member, the electric motor remains rotating due to the frictional load of the clutch.\nAnother method is to provide a limit switch which is actuated by the moveable member upon reaching the infinite focusing position to interrupt the flow of current to the motor. This, also has not been satisfactory. For one thing, it is difficult to completely accord an actuating position to the limit switch with the infinite focusing position of the objective lens. In the case where the switch is actuated before the focusing lens arrives at the stopper, defocusing results. If the switch actuating position is located slightly behind the mechanical stop position than all of the above mentioned problems, resulting from applying current to the motor after mechanical stoppage, will be incurred.\nAn appreciation of the criticallity in stop position can be had when it is realized that the allowable error of position of the objective focusing lens at infinite focusing is of the order of 0.01 mm. It is desirable to have a fine tuning or correcting means for adjusting the mechanical stop position so as to accurately position same. If a limit switch is used to stop the motor, the latter would have to be separately but relatively identically adjusted causing further difficulty."}
{"text": "1. Field of the Invention\nThis invention relates to an efficient method for neutralizing an alkoxylation catalyst. More particularly, this invention relates to a new use for partially spent ion exchange resins and an environmentally attractive method of neutralizing alkoxylation catalysts and removing potassium ions from polyol and alkoxylated alkyl phenol products. Example 2 demonstrates a reduction in potassium to less than 1 ppm.\nThis method is not only efficient and inexpensive, but also reduces problems inherent in disposal of a partially depleted acidic resin catalyst and reduces the impact on the environment of disposing of filter cakes which remain after using currently available methods to neutralize soluble base catalysts used in alkoxylation reactions.\n2. Description of Related Art\nOne of the most critical fields of study in industrial chemical processes is that of devising methods to remove impurities from products and effluent streams at low cost and, also, to avoid having to dispose of chemical compositions such as catalysts which are still partially active. Plants producing or using basic materials such as ammonia, organic amines, alkali salts or caustic have waste streams that often need to be neutralized. The manufacture of polyols and alkyl phenol alkoxylates presents such a problem.\nPolyols and alkyl phenol alkoxylates ar usually manufactured by the base catalyzed addition of propylene oxide or ethylene oxide to various initiators. The catalyst used is potassium hydroxide or another soluble base. The amount of potassium in the final product is critical in order to maintain control of downstream reactions. With one procedure, the catalyst is removed by neutralizing the base with a small excess of H.sub.2 SO.sub.4, absorbing the formed salt, (usually K.sub.2 SO.sub.4) with magnesium silicate, then filtering the product after heating to reduce the viscosity. The filter cake is then disposed of in a land fill.\nIn the past ion exchange resins have been used to produce polyols and, in fact, worked well. The disadvantage involved the fact that the solution used to regenerated the resin is considered a hazardous waste and is expensive to get rid of properly.\nThe use of ion exchange resins for removal of magnesium and noncarbonate minerals from water is known in the field of water treatment or water softening. The cost and frequency of regeneration are principal disadvantages. The ion-exchange material generally used in softening water are styrene-divinyl benzene copolymers. Kirk-Othmer Encyclopedia of Chemical Technology, 8, 70 and 24, 425 (1982).\nAnother aspect of the background for this invention is that the current transition to unleaded fuels in the U.S. has caused a demand for the addition of oxygenates into gasoline which has lead to the development of methyl-tert-butylether as a gasoline additive. The demand for MTBE has caused it to be the fastest growing chemical of the 80's and the demand will grow rapidly in the 90's. Chemical Business, January 1992, p.24. MTBE and other chemicals such as ethyl-tert-butylether (ETBE) and tert-amyl methylether (TAME) are made by the addition of alcohol to an olefin catalyzed by a sulfonic acid ion exchange resin.\nThe consumption of MTBE is currently about 180,000 barrels per day and could be as high as 670,000 barrels per day by the year 2000. Ibid, p. 25. Another reference indicates the demand would reach 1.2 million barrels per day by the year 2000. (Chemical Week, Nov. 20, 1991, p. 36.)\nDue to the demand for high productivity the ion exchange catalyst is used until reduced activity makes it more reasonable to exchange it for a fresh catalyst. This corresponds to a point of about 70-90%, and generally about 85%, of its original activity. This means there are enormous amounts of used catalyst that must be disposed of. Disposal of this partially spent catalyst gives rise to environmental concerns and often requires some sort of permit. It would be extremely efficient from a commercial and environmental viewpoint to accomplish a desirable goal using this partially spent catalyst, reduce the remaining activity of the resin so that disposal of it is not as objectionable and eliminate the production of a filter cake which requires disposal.\nIt would be a distinct advance in the art if an inexpensive means were available for neutralizing an alkoxylation catalyst and reducing impermissible levels of potassium in polyol and alkoxylated alkyl phenol products. It would be advantageous if this could be accomplished with a substance which to this point has had to be discarded while activity remains due to purity requirements of the primary product, i.e. MTBE, ETBE or TAME. To remove potassium from polyols and alkoxylated alkyl phenols with a partially spent catalyst, that normally entails disposal problems, and eliminate the neutralization of potassium hydroxide and the resulting filter cake would be quite efficient in every respect."}
{"text": "Virtual reality is a computer-generated simulation of an environment (e.g., a 3D environment) that users can interact with in a seemingly real or physical way. A virtual reality system, which may be a single device or a group of devices, may generate this simulation for display to a user, for example, on a virtual reality headset or some other display device. The simulation may include images, sounds, haptic feedback, and/or other sensations to imitate a real or imaginary environment. As virtual reality becomes more and more prominent, its range of useful applications is rapidly broadening. The most common applications of virtual reality involve games or other interactive content, but other applications such as the viewing of visual media items (e.g., photos, videos) for entertainment or training purposes are close behind. The feasibility of using virtual reality to simulate real-life conversations and other user interactions is also being explored."}
{"text": "The invention relates to a rotary bit for use in drilling boreholes or wells in underground formations. In particular, the invention relates to a rotary bit provided for cutting elements and a plurality of nozzles for discharging liquid, at least some of the nozzles being arranged to have liquid under pressure supplied thereto, thereby forming pressurized liquid jets. These jets either break up the bottom of the hole to deepen the hole (which manner of drilling is often indicated by the expression \"hydraulic drilling\") or remove the drilling flour or fine earth particles from the cutting elements and/or from the bottom of the hole, which drilling flour results from the mechanical drilling action of cutting elements carried by the bit. Such cutting elements may be mounted on roller cones and consists of cutting teeth, cutting rings, etc., or be mounted directly on the body of the bit and consist of diamonds, abrasive bodies such as the bodies made of materials like those known by the tradename Stratapax, and like elements.\nThe most effective action of the liquid jets will be obtained by jets originating from liquid nozzles having a relatively small internal diameter (in the order of 2-4 millimeter), over which nozzles a relatively high-fluid pressure difference exists (in the order of 30-150 bar). It will be appreciated that such nozzles that are made of an erosion-resistant material, will be liable to become plugged by particles present in the drilling mud that is being pumped down through the drilling equipment (such as the drill string) to the drill bit.\nThese particles may be lumps of solid material that has been mixed in powder form at the surface with a liquid for making up the drilling mud. Incomplete mixing will result in the formation of lumps that may partly be broken up by the mud pumps via which the mud is passed down the well, but the remaining lumps will be sieved off from the mud by jet nozzles in the bit, which nozzles will be plugged thereby decreasing the mud flow through the bit and jeopardizing the drilling action. Other particles may be constituted by fragments of the formation, which fragments have passed through damaged parts of the screening trays on the drilling floor, via which trays the drilling mud is returned to the borehole after appropriate treatment thereof. Also, dirt or corrosion products may be detached from the inner wall of the drill string by the mud flow passing therethrough, and be caught at the entrances of the nozzles when the mud passes therethrough.\nFurther, lost circulation material may be added to the mud flow to fight circulation losses occurring during drilling. These materials sometimes contain lumps or chunks of solid material that cannot pass through the small passages of the fluid nozzles present in the drilling bit.\nClogging of the fluid nozzles of the drilling bit described in U.S. Pat. No. 3,175,629 to D. S. Rowley is prevented by a screening element mounted in the shank of the bit, which screening element prevents relatively large particles present in the mud stream to reach the three fluid nozzles that debouch in the face of the bit. Only part of the mud stream is passed through the screening element, whereas the remaining part passes through a central channel in the shank and the body of the bit to a large-diameter nozzle situated in the centre of the lower end of the bit body. A plurality of choke plates provided with a single central passage of a diameter equal to the diameter of the central nozzle is arranged in the central channel to restrict the flow through the large-diameter central channel in order that a sufficient volume of drilling fluid will pass through the screening element to the small-diameter nozzles. The screening element is self-cleaning, and the particles caught by the screening element are subsequently discharged by the fluid flow passing through the central channel.\nA similar arrangement, apart from the choke plates in the central channel, is known from the U.S. Pat. No. 2,293,259 to C. D. Johnson. Part of the mud flow has the relatively large-sized particles filtered therefrom, and is supplied to six nozzles, whereas the remaining unfiltered part of the mud flow is passed through a nozzle of a diameter larger than the diameters of the six nozzles.\nTests carried out with drilling bits have made clear that drilling efficiency will considerably be promoted by the use of liquid jets of extremely high pressure and high velocity. Unfortunately, when raising the pressure of the mud supply to the jet nozzles in the abovementioned prior art bits, large volumes of fluid will pass through the central channel of the bits to the centrally arranged large-sized nozzle without effectively supporting the drilling action of the bit. In the bit of U.S. Pat. No. 3,175,629, this situation might be improved by increasing the number of choke plates in the central channel, but this will lead to a complicated and costly structure of the bit."}
{"text": "1. Field of the Invention\nThe present invention relates to an electronic apparatus provided with facsimile transmission means for facsimile transmission of two-side image data, and a control method and a control program therefor.\n2. Related Background Art\nThe conventional facsimile apparatus is basically designed on an operation of transmitting an original bearing information on one side thereof and recording such information on one side of a sheet. Also the conventional ITU-T Recommendation T4 or T30 describes such one-side communication only.\nOn the other hand, with recently increasing attention toward resource saving, so-called two-side recording mode for recording on both a front side and a reverse side of the recording sheet is increasingly provided in a printer, a copying machine, etc. Also the standard for the two-side facsimile communication was recommended by the ITU-T in February 2000. This is the ITU-T Recommendation T.30 which defines informing the presence or absence of a receiving function for the two-side information from a receiver to a transmitter and informing whether the communication is in the two-side mode from the transmitter to the receiver.\nAlso for transmitting the two-side image information, there are known two-side alternate transmission (alternate mode) of transmitting an image on the front side of an original and that on the reverse side alternately for each page, and two-side continuous transmission (continuous mode) of transmitting all the images on the front sides of the originals and then transmitting all the images on the reverse sides.\nOn the other hand, in a facsimile apparatus executing the conventional one-side communication, there is known an apparatus capable of so-called collective transmission.\nIn such collective transmission, in case plural image data to be transmitted to a same address are present in a transmission image memory on account of various reasons, for example, the postponement of the transmission owing to the destination being busy at the timer transmission or at the previous calling, the image data to the same address are collectively transmitted in one call connection.\nSuch collective transmission is effective in reducing the communication charge and is therefore preferably utilizable also in the facsimile apparatus capable of two-side communication, but such collective communication is not considered in the present Recommendation for two-side communication.\nFor example, since the aforementioned Recommendation T.30 assumes that all the original images to be transmitted in one two-side communication are two-side images, a machine constituted simply based on such Recommendation cannot achieve the collective communication in case a one-side original and a two-side original are mixedly present for the same address. On the other hand, a configuration simply designed to execute collective communication by an independent procedure may hinder interchangeability with a facsimile apparatus made by another manufacturer."}
{"text": "1. Field of the Invention\nThis invention relates to digital transmission line termination circuits, and more particularly, to a termination network which combines the functionality of a terminating impedance with protection to the digital transmitter/receiver circuitry, protection from faulty differential receiver output due to its loss of input, and compatibility with various differential receiver technologies.\n2. Description of the Prior Art\nIt is known that termination networks are necessary for digital circuits operating at very high frequencies. High speed logic circuits can produce digital pulses with logic level transitions of very short duration, which in turn can cause reflections or \"ringing\" on the signal path if these transition times approach the propagation delay of the transmission line. To avoid these effects, termination resistors or termination networks which match the characteristic impedance of the transmission line are placed in parallel with the receiver circuit in order to dissipate the energy of the pulse and prevent reflections or \"ringing\".\nIt is also known that \"differential\" transmitter and receiver circuits have been used to eliminate noise problems on high speed transmitter/receiver interfaces, or on transmitter/receiver interfaces with long transmission lines. A differential transmitter will simultaneously send both a signal and its inverse to the differential receiver, rather than sending only a single signal. Ground voltage differences are therefore nullified between the transmitter circuit and receiver circuit, which provides high noise immunity. The use of differential signals also avoids the problem of voltage drop on the transmission line that is associated with \"single-ended\" transmission lines. A single-ended transmission line experiences a voltage drop due to the resistance of the transmission line, and this voltage drop decreases the low voltage threshold at the receiver circuit. Differential transmission signals have this voltage drop on both signals, whereby the effect of the voltage drop is cancelled out. A differential approach typically provides more than twice the noise margin of the corresponding single-ended approach.\nIn order to prevent reflections and ringing on a differential transmission line, the proper characteristic impedance of the differential transmission line must be maintained through the use of a differential termination network. A differential termination network must provide a termination impedance equal to the characteristic impedance of the transmission line at each of the two transmission line signals. It is well known that both signal paths of a high frequency differential transmission must be properly terminated in order to avoid reflections or ringing. There are inherent problems with prior known differential termination networks that occur in a system which requires the differential transmitter circuit to be powered by a voltage source separate from that of the differential receiver circuit and differential termination network.\nA differential receiver circuit is designed so that its output signal switches from one logic state to the other when its input signals approach equal voltage potentials. This can create a problem if the transmission line driving the receiver becomes open circuited, which can occur if the transmission line itself is removed or damaged, or if the voltage which supplies the transmitter circuit fails when the receiver circuit is powered by a separate voltage source than that of the transmitter circuit. In a system which utilizes two separate and independent power domains, it is not unlikely that the voltage supplying power to the transmitter can fail, while the receiver remains powered up. The problem which could occur in this instance is that the two inputs of the differential receiver could float to an equal voltage level, thus causing the receiver output to come to rest at an indeterminate voltage level, or to continuously toggle.\nThe present invention solves the problem at the differential receiver circuit when the differential transmitter circuit signals are not present. In order to ensure that the receiver circuit's output signal does not continuously toggle, a voltage differential must be maintained at the inputs of the receiver circuit upon loss of the transmitter circuit signals. U.S. Pat. No. 4,567,384, by Stuhlmiller, issued Jan. 28, 1986, shows an example of a circuit which may be used with a EIA RS 422/423 receiver to hold the output signal of the receiver to a known output level when the input signals to the receiver become unavailable. The Stuhlmiller design utilizes transistors, resistors, and diodes to create the offset voltage, and it is used in addition to termination impedances. The present invention combines the functionality of a terminating impedance with the ability to provide a voltage differential to hold the receiver circuit in an inactive state upon loss of the transmitter circuit signals, without the need for additional components.\nAnother problem which has prompted a design feature in this termination network is that different types of differential receiver circuits may require different minimum voltages at the inputs of the receiver circuit. For instance, one type ECL differential receiver may require 2.0 volts minimum at the inputs for the low threshold voltage, while another type may require 2.6 volts minimum at the inputs. As previously stated, if the transmitter circuit fails or the transmitter supply voltage fails, the two input signals to the receiver circuit must be offset so that the it's output does not toggle. Therefore, the two voltages need to be offset from one another, but also may need to maintain a certain minimum voltage in order for the differential receiver circuit to operate correctly. For instance, the differential receiver circuit may require 2.6 volts minimum at the inputs, but the voltage divider termination resistances may only produce 2.2 volts (ignoring the offset voltages). This problem can be solved by providing a threshold voltage within the termination network which increases the two differential receiver input voltages to a minimum threshold voltage. This voltage can then be adjusted according to the minimum required input voltage of the particular differential receiver circuit, without the need to change the resistance values of the termination resistors.\nThe present invention also embodies a feature which solves a problem where the voltage that powers the differential receiver circuit is separate and independent from the voltage source that powers the differential transmitter circuit. If the voltage source that powers the receiver circuit was to fail, the failed voltage source would appear as a system ground. Since the differential transmitter circuit is expecting the differential receiver voltage to be present, the transmitter circuit could overdrive its DC output current capability when the failed voltage source appears as a system ground. In order to prevent current overdrive of the transmitter circuit upon loss of voltage to the receiver circuit, a resistance is placed between the termination impedance and ground. This separate resistance will effectively raise the equivalent resistance in the termination network when the receiver circuit's voltage has disappeared. The required termination resistances will not be affected by this extra resistance during normal operation however, because a bypass capacitor is used in parallel with the resistor. Therefore, the transmitter circuit is protected when the receiver circuit voltage fails, but during nominal operation the protecting circuitry will be transparent to the rest of the system.\nOften Emitter Coupled Logic (ECL) is used for differential transmission because of its high speed characteristics. Although ECL circuitry has the favorable quality of high speed, ECL circuitry consumes a great deal of power. The resistance and bypass capacitor (which protects the transmitter circuit in the event of the receiver circuit's voltage failing) also operates to reduce the overall power consumption of the system during normal operation.\nThe present invention is intended to alleviate problems where a transmitter circuit and an associated receiver circuit are powered by different voltage sources. The termination network also provides current limiting to reduce power consumption, and provides flexibility in choosing a differential receiver circuit."}
{"text": "Content may be added to and/or overlaid over a document (i.e., an image, a web page, a video, etc.) before the document is provided to a user. For example, messages, text corresponding to dialog (e.g., closed captioning), internet data (e.g., search results), messages from social media, data from another document, advertisements, etc., may be added to a document. Attributes of the additional content, such as a color of the additional content, is selected regardless of displayed attributes of the document."}
{"text": "How to effectively control and mitigate the pollution and hazard of heavy metals in soils is an increasingly serious international problem, which is especially prominent in our country. In particular, it is difficult to find a low-cost, broad-spectrum and practical technical method for treating large-scale heavy-metal polluted farmland soils. Besides, it is especially difficult for treating heavy metal composite pollutants in farmland. For example, arsenic and cadmium composite pollutants, compared to a single pollutant, result in more complex environmental effects due to their interaction, and are more difficult to be treated. Increasing the pH of soil can effectively reduce the bioavailability of cadmium in soil and the content of cadmium in rice, but may increase the activity of As in soil. Moreover, flooding can reduce the content of Cd in rice, but increase the content of As in rice. Therefore, how to control the soil process and manage water and fertilizer under the condition of As/Cd composite pollution is very complicated.\nIn recent years, the technology of passivating heavy metal pollutants in farmland has received more and more attention. A soil-friendly conditioner is applied to passivate heavy metals in polluted soil and reduce the absorption of heavy metals by crops, thus realizing simultaneous control and use of the agricultural soil polluted by heavy metals, which is a new economical and safe idea of prevention and control of heavy metal pollution in soil. The soil heavy-metal passivation technology has relatively low investment, high repair efficiency, and simple operation; it has superiority for repairing soils polluted by heavy metals at a medium/low level, and can meet the current demand of our country for controlling heavy metal pollution of farmland soil and ensuring safety of agricultural products. Commonly used soil heavy-metal passivators include the following substances: lime, calcium carbonate, fly ash and other alkaline substances; hydroxyapatite, ground phosphorite, calcium hydrogen phosphate and other phosphates; natural or modified zeolite, bentonite and other minerals; blast furnace slag, slag and other silicon-containing fertilizers; and peat, farmyard manure, green manure, bio-organic carbon and other organic fertilizers. These passivators often have a good passivation effect on heavy metal pollution in soil. However, the heavy metal pollution of soil is often composite pollution resulted from coexistence of two or more kinds of metals. Different heavy metal ions have different physical and chemical properties, and their migration in the soil and their environmental behavior are also quite different. Using a single soil passivator to repair multi-metal polluted soils is often challenging, and it is difficult to find a single substance that can reduce the mobility of all heavy metal ions. At present, therefore, the passivation of soil is mainly focused on passivation of a single heavy metal element, and composite additives or a variety of repair methods are usually simultaneously used for multi-metal composite pollution. If a multi-functional heavy-metal passivation material that can passivate multiple pollutants at the same time can be prepared, it will inevitably reduce the repair cost and improve the repair efficiency.\nThe toxicity and bioavailability of heavy metals in soil are not only related to the total amount of heavy metals, but also mainly affected by the physical and chemical properties of the soil. Iron oxides in soil are the key factors that control the morphological transformation and bioavailability of heavy metals in the soil. Iron is the most important redox active element in red soils (Wang et al., 2009) with a geochemical abundance of 5.1%, ranking fourth (Zhao Qiguo, 2002); iron is trapped in the surface of soil particles mainly in the form of free iron oxides, having high geochemical activity; it directly affects many soil processes (Borch et al., 2010). Iron oxides and other iron minerals, due to their very large specific surface area, chemical activity and morphological transformation, have higher adsorption capacity for a large number of heavy metals and oxygen-containing anions (such as PO43−, AsO43−, CrO42−, etc.), and are often used as treatment agents for wastewater polluted by arsenic, chromium and other heavy metals. However, because iron and its oxides are easy to aggregate, they combine with the elements Mg, P, Ca and S in soils and affect their chemical properties, which hinders the application of iron and its oxides in the repair of heavy-metal polluted soils.\nSulfur in soils plays an important role in controlling the activity and bioavailability of heavy metals. SO42− entering the soil is quickly reduced to S2− under anaerobic conditions, with S2− forming sulfides with metal ions to stabilize heavy metals; SO42− generated by organosulfur mineralization of the soil and SO42− entering the soil due to atmospheric deposition and fertilization are rapidly reduced to S2− under anaerobic conditions, with S2− forming sulfides with metal ions. Metal sulfides in anaerobic soils are stable and insoluble, and have a significant impact on the concentration of heavy metal ions in soil pore water. In paddy fields, therefore, the application of sulfur-containing fertilizers, especially at the rice seedling stage, plays a very important role in stabilizing heavy metals in soils. However, when sulfur is oxidized in the soil to form SO42−, a lot of H+ is generated, resulting in the activation of heavy metals. Therefore, the application of sulfur-containing fertilizers alone has the risk of reactivating heavy metals in the later stage of rice growth. If the sulfur-containing fertilizers are applied together with other soil conditioners to avoid oxidation again, the passivation effect and application scope of sulfur-containing passivators for heavy metal cadmium will be inevitably increased.\nSilicate fertilizer can reduce the availability of heavy metals in soil and inhibit the absorption and accumulation of heavy metals by crops. In recent years, studies have shown that silicon can increase resistance of plants to heavy metal toxicity, and is easy to use and cheap, thus having attracted people's attention. The current research has shown that the application of silicon fertilizer can increase the resistance of plants such as rice to heavy metals such as manganese, iron, cadmium and aluminum, and reduce the absorption and accumulation of these heavy metals by plants such as rice. Although the total content of Si in soil is very high, Si exists in the form of silicate; the silicon that plant can absorb and use is monosilicic acid (Si(OH)4), and so the content of effective silicon in soil is often very low. The content of silicic acid in soil solution is typically 0.1-0.6 mM (Epstein, 1994). With the development of intensified agriculture, long-term continuous cropping of crops (especially gramineous plants such as rice) will result in continuous absorption of the effective Si from the soil, leading to a reduction of crops. Due to the strong desiliconization in the red soil area, the content of the effective silicon in the soil tends to be lower. Therefore, the role of silicon fertilizer in agricultural production has drawn more and more attention. However, most of the current silicon fertilizers come from blast furnace slag, silicon-bearing ores and the like, and the silicon in these silicon fertilizers is less effective. Silicic acid and other silicon-containing salts, once applied to the soil, are also easily fixed by soil minerals.\nBiochar is a new type of material, and is a product of pyrolysis of biomass under hypoxic conditions. Biochar, whose particles are small, finely distributed, and light in weight, is a black porous solid, and is mainly composed of carbon, oxygen and other elements, with the carbon content therein generally above 70%. Biochar can be prepared from a wide range of raw materials, such as sawdust, straw, industrial organic waste, municipal sludge, palm filaments, coconut filaments, etc. Biochar is characterized by loose porous structure, large specific surface area, a lot of negative charge on the surface and so on, which make biochar have good adsorption properties. As a soil structure improver or a repair agent for soil pollution, biochar can improve the pH value of acidic soil and increase the cation exchange capacity, thus adsorbing the pollutants and heavy metals in the soil, reducing the accumulation of Cd, Pb and Zn in the crop body, accelerating the microbial metabolism and improving the microbial biomass of the soil, thereby enhancing the soil fertility, increasing the output of rice and other crops, and improving the quality of agricultural products. However, mobility of As increases with the pH of soil and As tends to bond to oxides and hydroxides of Fe, Al and Mn that have anion exchange sites in the soil, which mean that adding biochar to the soil does not necessarily control the bioavailability of As. Zheng et al. found that when biochar was applied to paddy soil polluted by heavy metals, biochar promoted the formation of iron film in paddy soil, which affected the migration of Cd, Zn, Pb and As in the soil, with the concentrations of Cd, Zn and Pb in rice roots decreased by 98%, 83% and 72%, respectively, while the concentration of As increased by 327%. Most farmland soil pollution is heavy-metal composite pollution, which poses a challenge to passivation of heavy metals in farmland by biochar.\nThe above iron oxides, sulfates, silicates and biochar have been widely used in heavy metal passivation of soil. They each have a good passivation effect on certain heavy-metal polluted soil under specific soil conditions. For example, iron oxides have a good effect on the passivation of arsenic-polluted soil, biochar has a special effect on the passivation of cadmium-polluted acid soil, and sulfates and silicates have good passivation effects on Cd or Pb pollution of acid soil. However, under the condition of heavy metal composite pollution of soils, the application of one of the above-mentioned conditioners alone often cannot achieve the goal of simultaneously passivating various heavy metals, and two or more passivators need to be applied at the same time. This causes not only inconvenience to the application, but also the phenomenon that iron oxides, sulfates and silicates tend to react chemically, resulting in loss or decrease of the passivation effect."}
{"text": "1. Field of the Invention\nThe present invention is generally directed toward compositions for and methods of digitally printing an ink image onto a woven textile material, especially those comprising synthetic resin fibers. In particular, digital printing systems are disclosed that utilize a base application comprising an acrylic latex material that is formulated to be imprinted upon with digital ink while still wet, that is, before being fully cured. Accordingly, the base application permits rapid printing of an ink image onto the textile material without an intermediate drying or baking cycle in which the base application is fully dried or cured prior to deposit of the ink image layer.\n2. Description of the Prior Art\nWoven textiles made from synthetic resin materials have been utilized for sometime in the manufacture of articles of clothing. However, advances in synthetic fiber technologies have resulted in the ability to manufacture clothing, and shirts in particular, having enhanced performance characteristics, such as in wicking moisture away from the wearer's body. These improved performance characteristics have increased synthetic resin textile popularity in many areas, especially athletic wear.\nThe synthetic resin textile materials often contain dyes that impart a desired color to the fabric. It is also quite common for the clothing articles formed from the dyed textiles to bear imprinted images, such as artwork, text, logos, etc., that customize the article based upon the user's preferences. Screen-printing of these various images is a popular printing method used to provide this customization of clothing articles, especially T-shirts. However, screen-printing of images has certain drawbacks. Preparation of the screens required for the printing process can be quite expensive depending upon the complexity of the image to be printed. Thus, in order to be economically feasible, large quantities of the particular clothing article must be prepared. Accordingly, screen-printing is generally not ideal for single or small production runs. Also, inks used in screen printing operations generally are more viscous and must be deposited onto the textile material in relatively thick layers that can negatively affect the feel of the clothing article as it is being worn.\nIn view of these drawbacks, digital printing has been explored as an alternative to screen-printing for creation of images on woven textile materials. However, present digital printing technologies for textile materials, especially synthetic resin textile materials, also exhibit certain shortcomings. The inks used in digital printing are generally less viscous than screen-printing inks and may tend to penetrate more deeply into the textile fibers thereby degrading image quality. In order to prevent the inks from absorbing too deeply into the fibers, curing of the resin systems making up the inks may be accelerated by exposing the textile material to elevated temperatures. Commonly, the textile material is passed through a dryer operating at a temperature of between 285-300° F. immediately after the image is imprinted. However, exposure to these elevated temperatures, which is commonly for 6-8 minutes, can deactivate the textile dye system causing the dye to weep into and degrade the printed image. Particularly in the context of synthetic resin textile materials, long-term adhesion of the cured ink image to the textile fibers can also be problematic. If the adherence of the ink to the textile fibers is weak, the clothing article will exhibit poor washability as the ink image may crack, fade, or release upon undergoing plural wash/dry cycles.\nImage transfer systems, wherein the ink image is first imprinted upon a temporary carrier substrate and then subsequently transferred to the clothing article, have also been proposed. But, these systems are generally more labor and material-intensive, thereby adding expense to the overall cost of producing the clothing article. Accordingly, there is a real and unfulfilled need in the art for a system of printing an image onto a textile material that avoids the aforementioned problems with existing screen-printing and digital printing technologies."}
{"text": "1. Field of the Invention\nThe invention relates generally to the field of watercraft. More particularly, it relates to apparatuses and methods useful for helping to center watercraft—such as personal watercraft, ski boats, fishing boats, luxury boats, and the like—as they are docked (such as on trailers or watercraft-lifting systems).\n2. Description of Related Art\nSome trailers for loading and transporting boats and other watercraft are equipped with bunks positioned along the length of a portion of the trailer. As the boat or other watercraft is driven onto the trailer for loading, the hull of the craft contacts the bunks and ultimately rests on them once the craft is out of the water. The bunks help guide a craft that is loaded in calm waters, but they are not very useful for centering a craft that is loaded in rough waters and/or strong crosswinds. A craft loaded during such conditions can get banged around by the bunks and potentially damaged.\nSome trailers use rollers (e.g., wobble rollers and keel rollers) instead of, or in addition to, bunks. Rollers are not much more helpful than bunks at centering a craft that is loaded in rough waters and/or strong crosswinds. Furthermore, rollers may damage the hull of the craft because of the large amount of force that can apply to a small area of the hull.\nA number of trailer guides and mechanisms designed to make boat loading easier have been disclosed. See U.S. Pat. Nos. 4,209,279; 3,608,754; 3,026,981; 3,447,815; 5,228,713; 5,013,206; 4,242,768; 4,715,768; 4,094,527; 4,099,279; 5,299,903; 4,500,249; and 4,529,217, all of which are incorporated by reference. While each of these disclosures purport to solve problems associated with boat loading, the inventor has discovered that none is completely satisfactory.\nSome watercraft, especially boats, are stored on the water instead of being trailered. Watercraft that are stored on the water are generally lifted out of the water using a lifting mechanism of some kind in order to minimize damage to the craft that might otherwise occur (i.e., corrosion from the water, damage due to rough waters, etc.). Some lifting mechanisms involve lift tanks to which bunks are attached. Although lift tanks are generally more protected from rough waters and/or crosswinds than are trailers, such conditions can still make steering a craft into place over a lift tank (or other watercraft-lifting system) difficult. The bunks attached to such lift tanks are generally not satisfactory at helping to center or centering the craft in the appropriate position over the tanks."}
{"text": "Existing electronic equipment cabinets generally employ forced air convection cooling of electronic modules mounted in chassis within the cabinet. Many prior art cabinet and chassis configurations circulate cooling air in a side-to-side configuration. The side-to-side cooling pattern has the disadvantage of pre-heating air for each successive chassis when chassis are placed side-by-side as is common in telecommunications and other embedded applications. As processor speeds increase and electronic modules generate more thermal energy, prior art configurations and methods of cooling electronic modules and their cabinets will become inadequate.\nAccordingly, there is a significant need for an apparatus and method that overcomes the deficiencies of the prior art outlined above.\nIt will be appreciated that for simplicity and clarity of illustration, elements shown in the drawing have not necessarily been drawn to scale. For example, the dimensions of some of the elements are exaggerated relative to each other. Further, where considered appropriate, reference numerals have been repeated among the Figures to indicate corresponding elements."}
{"text": "1. Field of the Invention\nThe present disclosure relates to image processing, and more particularly to a system and method for estimating a size of solitary objects in of interest in medical images.\n2. Description of Related Art\nMeasuring the size of pulmonary nodules from X-ray computed tomography (CT) data is an important practice for diagnosis and progression analysis of lung cancer. The nodule size often plays an important role in choosing a proper patient care, and is also an effective feature to separate true nodules form nodule-like spurious findings. Typically, the size is represented by the diameter of the nodule. Automating this task for computer-aided diagnosis (CAD) is, however, a difficult problem due to intensity variations, partial volume effects, attachment to other structures, and noises.\nA Ct-based screening protocol specified by the International Early Lung Cancer Action Program (I-ELCAP) details how the diameter of pulmonary nodules should be measured and how the measurements should be used for determining the patient management. According to the protocol, the result of an initial CT screening of lung is considered positive if at least one solid or part-solid nodule with 5.00 mm or more in diameter or at least one non-solid nodule with 8.0 mm or more in diameter is found. Although these 5 mm and 8 mm thresholds are likely to drop as more accurate screening becomes possible with high resolution multi-detector helical CT (MDCT), the importance of module size in cancer diagnosis will stay unchanged.\nTo this end, an automated size estimation algorithm using a Gaussian Ellipsoid Fit (EF) has been developed. Given a marker positioned near a nodule, the algorithm computes the location, orientation and radii of an ellipsoid that models closely the intensity variation nearby the marker. It employs mean-shift and scale estimator to find the solution. The volume and diameter of the nodule can be estimated from the ellipsoid. EF can be incorporated into a CAD system where markers and provided from either manual markings of a reader or the detection algorithm of the system. The accuracy of the estimates has been verified using a large database with manual size measurements. However, the technique tends to have difficulties with small nodules due mainly to the small sample size problem."}
{"text": "This invention relates to tool holders, specifically to holders which are used for holding pneumatic and electric hand tools when working on ladders, scaffolding and other precarious elevated working environments."}
{"text": "This invention relates to N-substituted benzylic aniline derivatives that are inhibitors of mitotic kinesins (in particular the mitotic kinesin KSP) and are useful in the treatment of cellular proliferative diseases, for example cancer, hyperplasias, restenosis, cardiac hypertrophy, immune disorders and inflammation.\nAmong the therapeutic agents used to treat cancer are the taxanes and vinca alkaloids. Taxanes and vinca alkaloids act on microtubules, which are present in a variety of cellular structures. Microtubules are the primary structural element of the mitotic spindle. The mitotic spindle is responsible for distribution of replicate copies of the genome to each of the two daughter cells that result from cell division. It is presumed that disruption of the mitotic spindle by these drugs results in inhibition of cancer cell division, and induction of cancer cell death. However, microtubules form other types of cellular structures, including tracks for intracellular transport in nerve processes. Because these agents do not specifically target mitotic spindles, they have side effects that limit their usefulness.\nImprovements in the specificity of agents used to treat cancer is of considerable interest because of the therapeutic benefits which would be realized if the side effects associated with the administration of these agents could be reduced. Traditionally, dramatic improvements in the treatment of cancer are associated with identification of therapeutic agents acting through novel mechanisms. Examples of this include not only the taxanes, but also the camptothecin class of topoisomerase I inhibitors. From both of these perspectives, mitotic kinesins are attractive targets for new anti-cancer agents.\nMitotic kinesins are enzymes essential for assembly and function of the mitotic spindle, but are not generally part of other microtubule structures, such as in nerve processes. Mitotic kinesins play essential roles during all phases of mitosis. These enzymes are “molecular motors” that transform energy released by hydrolysis of ATP into mechanical force that drives the directional movement of cellular cargoes along microtubules. The catalytic domain sufficient for this task is a compact structure of approximately 340 amino acids. During mitosis, kinesins organize microtubules into the bipolar structure that is the mitotic spindle. Kinesins mediate movement of chromosomes along spindle microtubules, as well as structural changes in the mitotic spindle associated with specific phases of mitosis. Experimental perturbation of mitotic kinesin function causes malformation or dysfunction of the mitotic spindle, frequently resulting in cell cycle arrest and cell death.\nAmong the mitotic kinesins that have been identified is KSP. KSP belongs to an evolutionarily conserved kinesin subfamily of plus end-directed microtubule motors that assemble into bipolar homotetramers consisting of antiparallel homodimers. During mitosis KSP associates with microtubules of the mitotic spindle. Microinjection of antibodies directed against KSP into human cells prevents spindle pole separation during prometaphase, giving rise to monopolar spindles and causing mitotic arrest and induction of programmed cell death. KSP and related kinesins in other, non-human, organisms, bundle antiparallel microtubules and slide them relative to one another, thus forcing the two spindle poles apart. KSP may also mediate in anaphase B spindle elongation and focussing of microtubules at the spindle pole.\nHuman KSP (also termed HsEg5) has been described [Blangy, et al., Cell, 83:1159-69 (1995); Whitehead, et al., Arthritis Rheum., 39:1635-42 (1996); Galgio et al., J. Cell Biol., 135:339-414 (1996); Blangy, et al., J. Biol. Chem., 272:19418-24 (1997); Blangy, et al., Cell Motil Cytoskeleton, 40:174-82 (1998); Whitehead and Rattner, J. Cell Sci., 111:2551-61 (1998); Kaiser, et al., JBC 274:18925-31 (1999); GenBank accession numbers: X85137, NM004523 and U37426], and a fragment of the KSP gene (TRIP5) has been described [Lee, et al., Mol. Endocrinol., 9:243-54 (1995); GenBank accession number L40372]. Xenopus KSP homologs (Eg5), as well as Drosophila K-LP61 F/KRP 130 have been reported.\nCertain quinazolinones have recently been described as being inhibitors of KSP (PCT Publ. WO 01/30768, May 3, 2001).\nMitotic kinesins are attractive targets for the discovery and development of novel mitotic chemotherapeutics. Accordingly, it is an object of the present invention to provide compounds, methods and compositions useful in the inhibition of KSP, a mitotic kinesin."}
{"text": "The present invention relates to microprocessor controlled embedded systems and, in particular, to the safety critical monitoring thereof.\nMicroprocessor controlled embedded systems are prevalent in medical devices, such as MRI imaging systems, blood testing equipment, etc., where it is now necessary to provide a safety critical environment for the protection of patients. As shown in FIG. 1, in a typical embodiment the system 2 sits between the specific application system 1 and an application specific user interface 3. In systems of this type, as shown in more detail in FIG. 2, application specific program logic resides permanently within a static memory 8 (non-volatile ROM or RAM). Along with the application specific program logic, the system includes an embedded operating system which is either a single process executive or multi-tasking executive system. The operating system is responsible for controlling the system interrupts, input/output control and the scheduling of the specific application tasks on a time or performance basis. The operating system operates with the microprocessor 5 which is part of the central processing unit of the system along with a random access memory and an input/output circuit 7 which communicate on a bus 6.\nThe application specific program logic is responsible for performing all functional operations, as well as for directly handling input and output of the information and control required for the specific application system utilizing the embedded system. The operating system is a stand-alone program or set of programs with a specifically defined interface for connecting it to the application specific program logic. The application specific program logic has the required interface logic to link it with the operating system. In conventional embedded systems, the operating system has no knowledge of the performance requirements of the application specific program logic or the specific application system, but simply acts as the controller for the central processing unit."}
{"text": "1. Field of the Invention\nThe present invention relates to a wire bonding device, in particular, to a wire bonding device in which relative positional deviation of a capillary with respect to an object of bonding is caused by changes in temperature.\n2. Description of the Related Art\nTo assemble a semiconductor device, wiring is automatically performed on an electrode pad of a semiconductor substrate. For such wiring, an automatic bonding device is used. An automatic bonding device is equipped with a horn device for applying an ultrasonic wave to a bonding line paid out from the capillary. A bonding target point of a semiconductor device is monitored by a camera. Such an automatic bonding device is disclosed in Japanese Patent Application Laid-open No. Hei 01-161727. In such a well-known device, there are provided, in order to perform bonding positioning with high accuracy, a temperature detection means for detecting a temperature of a horn, and a correction means for obtaining a change in a length of the horn due to thermal expansion based upon a detection result of the temperature detection means to correct a previously set relative distance between a capillary and a camera.\nWhen coefficients of thermal expansion and dimensions of the support means for securing the camera in position and the horn and the relative positional deviation between the camera and the capillary from the initially set state due to changes in temperature are not taken into account, it is impossible to eliminate the positional deviation factors due to thermal expansion of the mechanism parts as a result of temperature change, and further, it is impossible to correct the deviation of the reference bonding position of the semiconductor device which is the object of bonding. The above-mentioned well-known technique does not teach any correction means for correcting thermal expansion of the support means for securing the camera in position, nor does it suggest the importance of such correction. Further, the thermal deformation of the object of bonding is not taken into account at all.\nIt is required that the absolute positional deviation of the camera should be taken into account and that the positional relationship between the camera and the horn should be corrected. For a future semiconductor device, there is also a requirement for a further improvement in the accuracy of the bonding position, which is to be attained by taking into account the thermal expansion/contraction of the object of bonding.\nIt is an object of the present invention to provide a wire bonding device in which the absolute positional deviation of the camera is taken into account and in which the positional relationship between the camera and the horn is corrected.\nIn accordance with the present invention, there is provided a wire bonding device comprising: an XY table (2), a first holding member (8) supported by the XY table (2), a camera (9) supported by the first holding member (8), a second holding member (5) supported by the XY table (2), a capillary (7) held by the second holding member (5), a stationary base (3) for securing an object of bonding (4) in position, and a driving amount calculating unit (13) for computing the driving amount of the XY table (2) on the basis of the positional deviation of the camera (9) with respect to the origin of the XY table (2) due to temperature change in the first holding member (8) and the positional deviation of the capillary (7) with respect to the origin of the XY table (2) due to temperature change in the second holding member (5).\nIn a coordinate system in which the camera (9) and the capillary (7) are caused to move integrally by the XY table (2), the first holding member (8) and the second holding member (5) undergo thermal deformation independently, with the result that a minute deviation is caused in the positional relationship between the camera (9) and the capillary (7), and this minute deviation is calculated, correcting the position of the object of bonding (4) detected and calculated through imaging of the camera (9). On the basis of the positional deviation of the capillary (7) held by the second holding member (5), the moving amount of the XY table (2) is corrected, whereby the deviation of the bonding position due to thermal fluctuations is eliminated, making it possible to realize a bonding of high positional accuracy. The relative positional deviation between the camera (9) and the capillary (7) due to thermal fluctuations can be physically calculated by detecting the temperatures of the holding member (8) supporting the camera (9) and the support member (5) supporting the capillary (7).\nThere is further added an expansion coefficient calculating unit for calculating the expansion coefficient of the object of bonding (4). In this case, the driving amount calculating unit (13) performs calculation by adding the positional deviation of the bonding point of the object of bonding (4) with respect to the stationary base (3) on the basis of the expansion coefficient, to both of the above-described positional deviations. By adding the positional deviation of the bonding (object) point due to the expansion of the object of bonding (4) to both of the above-described positional deviations, it is possible to perform bonding position control with higher accuracy.\nThe expansion coefficient can be calculated on the basis of a reference distance (L) between two points (24-1 and 24-2, or 25-1 and 25-2) of the object of bonding (4) and the present distance (L+xcex94L) between the two points at the present time (24xe2x80x2-1 and 24xe2x80x2-2, or 25xe2x80x2-1 and 25xe2x80x2-2). When the temperature is known, it is possible to calculate the positional deviation of a local point of the object of bonding on the basis of the known expansion coefficient of the object of bonding. Direct measurement of the expansion/contraction of the distance between the reference points of the object of bonding by the camera (9) makes the detection of the positional deviation amount more accurate.\nIt is possible to further provide a temperature detector (18) for detecting the temperature of the second holding member (5). In this case, the driving amount calculating unit (13) calculates the positional deviation of the capillary (7) based on the temperature. It is possible to provide a first temperature detector (16) for detecting a first temperature of the first holding member (8) and a second temperature detector (18) for detecting a second temperature of the second holding member (5). In this case, the driving amount calculating unit (13) calculates the positional deviation of the camera (9) on the basis of the first temperature, and calculates the positional deviation of the capillary (7) on the basis of the second temperature. In this case also, the expansion coefficient of the above-mentioned object of bonding (4) is taken into account."}
{"text": "Over the last fifty years, approaches toward providing animal nutrition have changed. No longer are the animals fed whatever grain or forage may be available. Instead, the diets of animals are closely monitored for total nutrition value, and for cost. The animal on the diet is monitored, for quality and performance characteristics, and for the environmental impact of the waste from the animal. The information gathered is employed to adjust the feed to increase nutrition value of the feed and the animal performance characteristics while decreasing the cost and environmental impact.\nCereals account for about half of all feed ingredients, primarily because they are good sources of energy. Maize tends to be the preferred feed grain because of its highly digestible carbohydrate and relatively low fiber content, which is particularly important for swine and poultry (Hard, Proc. Southwest Nutr. Conf. 43-54 (2005)). Because of the low protein content of maize, it is common practice to use feed additives and supplements, such as protein-rich feeds, amino acids, vitamins, minerals and fats in diets for swine and poultry. The ratio of cereals to supplements has changed through the years in an attempt to maximize feeding efficiency of the animals. The feeding efficiency (the feed conversion ratio or how much feed is required to produce one pound of animal weight) is determined by the genetic potential of the animal and by the nutrients supplied to the animal. As the feed conversion ratio has risen due to genetic enhancements, the mineral and nutrient requirements for feed necessary to assure a complete and healthy diet have risen. Since an animal's ability to feed limits the amount of nutrients and calories it can consume, the feed industry has had to develop ways to make feeds that have improved protein quality (improved balance of essential amino acids), digestibility (fiber, starch, anti-nutrients), and metabolizable energy (oil).\nSources of feed protein have come under global public scrutiny in recent years because of the bovine spongiform encephalopathy, or mad cow disease, crisis associated with the feeding of meat and bone meal as the primary protein source in animal diets in many parts of the world. Plant protein sources, especially soybean meal, a residual product of the oil extraction process from soybeans relatively high in protein, have become a dominant alternative protein supplement used in feed following bans on using meat and bone meal in many parts of the world.\nPlant protein sources, however, may lack sufficient levels of essential amino acids required for adequate animal health, growth and performance. Requirements vary depending on the species and age of the animal. For example, the order of the top three limiting amino acids in feed composed of corn and soybean meal is lysine, threonine, and tryptophan for swine and methionine, lysine, and threonine for poultry. (FAO Animal Production and Health Proceedings, Protein Sources for the Animal Feed Industry, xi-xxv, 161-183 (2004)). These limiting amino acids must be available at specific minimum levels in order for the animals to use dietary protein efficiently. (Johnson et al. “Identification of Valuable Corn Quality Traits for Livestock Feed”, Report from the Center for Crops Utilization Research, Iowa State University, 1-22 (1999)). Crude protein in feed ingredients is not totally digestible for any species, for example corn protein is approximately 84% digestible by poultry and 82% digestible by swine (Johnson et al. (1999)). To compensate for this inefficiency, feed often contains excess protein that then results in high nitrogen excretion. More stringent environmental regulations are being imposed because high nitrogen excretion poses serious concerns to human health through ammonia or nitrate/nitrite pollution in soil and water. One solution to the problems of nitrogen pollution associated with animal feeding is to decrease crude protein in feed and supplementing feed with amino acids. A one-percentage point reduction in crude protein content in feed can yield about eight to ten percent reduction in nitrogen excretion. (FAO Animal Production and Health Proceedings, Protein Sources for the Animal Feed Industry, 161-183 (2004)). Supplementing of feed with amino acids can provide the required nutrients while decreasing excessive crude protein and can provide limiting amino acids when they are not sufficiently available.\nAdditionally, animals lack the enzymes necessary to digest the non-starch based polysaccharides present in soybean meal, and corn/soybean feed mixtures resulting in high manure volume and environmental impact. Approximately 65 to 70% of the total phosphorous in cereal grains is organically bound in phytate phosphorous, which is relatively unavailable to poultry and swine because they lack the enzyme phytase required to digest phytate, thus requiring inorganic phosphorous supplements. (Knowlton, J. Anim. Sci. 82(E. Suppl.):E173-E195 (2004)). The undigested phytate passes through the digestive system and leads to excretion of excess nutrients resulting in high manure volume and high levels of phosphorus in manure. Manure containing nitrogen and phosphorous at levels in excess of crop requirements results in environmental contamination especially of water resources caused by runoff. Enzymes, such as phytase, are commonly added to feed to increase digestibility. The addition of phytase can reduce the level of phosphorus released in animal waste to about half the previous level. However, the cost of phytase is about three times the cost of the conventional inorganic phosphorous supplements usually added to feed. (“Enhanced Animal Feed Will Be A Boon For The Environment,” Economic Perspectives, Agricultural Biotechnology, An Electronic Journal of the U.S. Dept. of State, Vol. 8, No. 3, September 2003).\nEnd-users of feed corn include livestock producer-feeders, feed manufacturers, corn millers and processors. Whether for a livestock producer-feeder who mixes and prepares their own feed or for a feed manufacturer who supplies a variety of feed products including complete ration feeds or nutrient supplements, each of the various ingredients necessary to produce the right combination of nutrients (i.e. protein, amino acids, enzymes, etc.) will need to be transported from site of production and/or processing to the site of the end-user. The availability, price, and transportation requirements and costs of each component of a particular feed will vary from year to year and in different geographical regions. Feed is usually formulated to meet nutritional requirements at a minimum dietary cost. The feed industry balances rations to supply nutrients at the least cost. Because of the variability of the supply and cost of nutrients and additives, livestock feeders and feed manufacturers would value corn traits that create substitutability for more expensive feedstuffs or additives.\nBecause feed is around 60% of animal production costs, any savings can be considerable, especially in large operations. Nutritionally enhanced corn which can deliver higher levels of important nutrients and metabolizable energy, and/or enhanced digestibility and bioavailability of nutrients would provide the following benefits: reduced feed costs per unit weight gain or production of eggs or milk; reduced animal waste, particularly nitrogen and phosphorous; reduced veterinary costs and improved disease resistance; improved processing characteristics to make the feed; and improved quality (Johnson, et al. (1999)). Cost savings can be achieved by using nutritionally enhanced corn through, for example, reduced cost for needed supplements and synthetic additives, reduced transportation costs associated with the shipping of each additive and ingredients to produce the additives, reduced cost in mixing numerous additives during feed processing, and reduced costs associated with disposal of excess volume of manure.\nBoth traditional plant breeding and biotechnology techniques have been used to develop maize plants with desired traits such as low-phytate, high-lysine, or high-oil maize. For example, U.S. Pat. No. 5,723,730 describes an inbred corn line used to produce a hybrid with elevated percent oil and protein grain.\nExamples of grain-based feed that provide improved animal nutrition and can reduce environmental impact of animal production are described by Chang et al. in U.S. Pat. Nos. 7,087,261 and 6,774,288 and in U.S. Publ. No. 2005/0246791.\nThere remains a need to develop inbred parental maize lines that contribute these desirable traits to the hybrids in which they are used. These traits may also include resistance to diseases and insects, tolerance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination and stand establishment, growth rate, stalk strength, root strength, ear retention, maturity and plant and ear height, are important. Selection of germplasm that possess the desired traits is required to develop novel, desirable plant germplasm for plant breeding.\nMaize is an important and valuable field crop. Thus, a continuing goal of plant breeders is to develop stable, high yielding maize hybrids that are agronomically sound. The reasons for this goal are obviously to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the corn breeder must select and develop maize plants that have the traits that result in superior parental lines for producing hybrids and that provide end-user value."}
{"text": "A wide variety of communications services are available including, for example, call waiting, call forwarding, call blocking, do not disturb services, customized messaging services, communications circles, etc. Generally, the services are implemented for a particular customer based on profile information relating to the customer's preferences. For example, a customer may have call forwarding service implemented so all calls to his or her home telephone number during business hours are forwarded to a network voice mail service. As another example, a customer may have call blocking service implemented so calls received from a specific number during evening hours are blocked.\nA customer's preferences with respect to communications service may change from time to time. Referring to the examples above, the customer may decide to have calls that were previously forwarded to the network voice mail service forwarded instead to an office telephone. With respect to the call blocking service, the customer may decide to extend the call blocking service to block calls from another specific number. To accommodate the change in preferences, the customer's profile information relating to the communications service may need to be changed.\nGenerally, the service provider providing the service makes the change in the customer's profile information relating to the service. The service provider typically makes the change because the service provider delivers the service, and thus, controls the delivery of the service.\nTo make a change in a communications service, a customer notifies the service provider. The customer may notify the service provider in a number of different ways, which include calling a customer service number, or using the Internet to reach the service provider's web site and communicating the desired changes. Some service providers allow a customer to call a feature access code (FAC) and provide change instructions.\nThe necessity of having the customer contact the service provider and provide the change instructions significantly slows the desired change in the communications services. The additional necessity of having the service provider implement the change instructions further slows the desired change. Some customers may plan ahead or be patient so a delay of a desired change to communications services may not be important. Most customers, however, desire their changes to communications services to take effect as close to immediately as possible.\nTherefore, there is a need for faster ways of implementing a customer's desired changes to his or her communications services. There is also a need for more convenient ways of implementing a customer's desired changes. In addition, there is a need for faster and more convenient ways of implementing a customer's desired changes to his or her communications services without sacrificing qualities such as accuracy and thoroughness in the implementation of the changes."}
{"text": "With the steady growth in recent years of vehicle electronic systems, the number of sensors used in automobiles has also risen dramatically. Current vehicles can contain forty to fifty of such sensors. Among the types of sensors commonly provided are pressure sensors in the vehicle fuel systems, air bag sensors, over-current and over temperature-protection sensors, oxygen sensors for exhaust gas systems, catalyst temperature sensor, etc. With most if not all of such sensors, signal conditioning will be necessary for most sensor elements. Raw sensors exhibit some level of offset voltage in output signal level when no stimulus is applied to the sensor. Typically, a potentiometer or a trim resistor branch circuit that is laser trimmed for calibration is employed to calibrate the output signal from the sensor.\nFor example, in the case of exhaust sensors the trim resistor is used in the automotive oxygen sensor output to ECU circuitry connection system. Hitherto, this trim resistor has been provided as an integral part of the sensor wiring output connector assembly. In one commercially provided connector assembly six male blade terminals are provided and four of the terminals are crimped to wires coming from the sensor. The remaining two terminals are insert molded into the connector assembly and one end of each terminal is soldered to a contact pad on the trim resistor substrate. An opening is provided on the side of the connector assembly to install the trim resister and allow access for the laser trimming operation. The opening is then closed and sealed with a cover and gasket.\nIn this prior art construction currently in use the ceramic substrate and terminals are pre-molded parts that are inserted as a pre-molded part into the connector assembly mold and then secured therein as a mold-over part. Hence the connector assembly construction must be designed to match the connector assembly family currently used in the given wiring system of the vehicle design. This current practice is thus limited in flexibility and adaptability when it is desired to provide resistance trim capability as an add-on to exhaust system sensors in automotive vehicle electronic systems. Mold tooling costs are also a problem when system design changes are mandated.\nThere is thus a need to provide a trim resistor connector assembly construction that can be economically manufactured, assembled and installed to be operably coupled as an add-on to new as well as existing sensor wiring assemblies. In addition, there is the continuing need to improve exhaust sensor performance by enhancing calibration circuit construction. There is also the ongoing need to reduce the expense of the calibrating means for the sensor, and one that can be employed to calibrate the sensor prior to installing in the vehicle."}
{"text": "1. Field of the Invention\nThe invention relates to a method of detecting geological structures and to obtaining stereoscope views from aerial and satellite photographs. In particular the method of the invention concerns detecting geological structures by viewing either aerial photographs or satellite photographs to produce a stereo or 3D effect and to enhance that effect to enable the structures to be more readily detected.\n2. Description of Related Art\nAerial photography and in particular the advent in photo-interpretation occurred in the 1960's. Vertical aerial photography includes imaging in the visible as well as the infra red, radar and microwave bands of the electromagnetic spectrum. Both conventional aerial photographs and infra red photographs have specific applications for the science of geology. Aerial photography is confined to enable study in detail of local areas. For a more extensive study a wider photographic coverage is required. For example the study of sedimentary basins on a regional scale to a continental or national scale requires wider photographic coverage. The need for such regional and national scales of study has led geologists to use Landsat or other satellite photographs at scales from 1:1,000,000 to 1:100,000. The invention is concerned with both aerial and satellite photographs.\nIt is currently believed that stereovision or three dimensional images cannot be obtained from Landsat data. The current practice is for geologists to view standard processed landsat photographs with the naked eye and observe geological structures such as linears. Linears are structures detected involving meansurements in two dimensions only. The detection of linears from photographs provide limited interpretative value unless the structure is investigated and confirmed by geophysical, geomophological or geological data.\nDue to this restricted two dimensional observation of photographs information obtainable in this way has been of limited value.\nVertical aerial photography is well established as a geological mapping tool. Vertical aerial photography not only expedites field mapping but also may reveal information that cannot readily be obtained in any other way. Electromagnetic radiation includes ultraviolet radiation, visible light and infra red radiation (IR). Infra red aerial photography is useful in enhancing the contrast of the terrain and bodies of water are sharply contrasted in relation to land and vegetation types may also be contrasted from one another. Thermal imagery has importance from a geological point of view because differences between various rock types and other aspects can also be distinguished. Geological structures such as faults, fractures, joints, folds, crushed zones and others are detectable because of the relationship of the terrain to the rock type.\nInfra red aerial photograph employes a yellow filter over the camera to filter out blue light. However, not all of the blue light is eliminated.\nSatellite photography employs a sensor unaided by a filter of any kind and it is possible to record the full spectrum. This is limited by the sensitivity and spectral band of the sensor and the range of physical colour bands available to display the electromagnetic radiation recorded.\nThe Landsat satellite produces photographs called Landsat masters and these can be obtained in either colour or black and white. The photographs are derived from multi-spectral scanner data (MSS) which comprises radiation in four spectral bands:\nBand 4 visible green\nBand 5 visible red\nBand 6 non-visible infra red\nBand 7 non-visible infra red\nThe colour master is a composite of three spectral bands printed in registery through non-visible infra red, red and green filters on photographic colour negative material. These masters are referred to as \"first generation archival masters\" and are used to print second generation customer prints or transparencies. The masters are usually at a scale of 1:1,000,000 and cover an area of 185 kilometers square. The picture elements (or pixels) cover an area of approximately 80 square meters.\nIt is an object of the invention to provide a method of detecting geological structures from aerial and satellite photographs.\nIt is another object of the invention to provide a method of stereoscopically viewing two photographs."}
{"text": "Multi-layer piezoelectric elements mounted on the fuel injection apparatus of automobile engine or the like have been known. FIG. 5A is a perspective view showing the multi-layer piezoelectric element of the prior art. The multi-layer piezoelectric element 101 comprises a stack 109 made by stacking a plurality of piezoelectric layers 105 and a plurality of internal electrode layers (metal layers) 107 which are stacked alternately one on another, and a pair of external electrodes 111 formed on the side faces of the stack 109.\nThe internal electrode layers 107 are not formed over the entire principal surface of the piezoelectric layer 105, and have the so-called partial electrode structure. FIG. 5B is an exploded perspective view explanatory of the partial electrode structure showing a part of the multi-layer piezoelectric element 101 shown in FIG. 5A. As shown in FIGS. 5A and 5B, the internal electrode layers 107 are stacked so as to be exposed on either side face of the stack 109 alternately in every other layer. Accordingly, the plurality of internal electrode layers 107 are electrically connected to the pair of external electrodes 105 alternately. The stack 109 has inactive layers 113 stacked on both ends of the stack 109 in the stacking direction.\nIn the multi-layer piezoelectric element which has the partial electrode structure shown in FIGS. 5A and 5B, there are an active region A in which the internal electrode layers 107 of different polarities oppose each other via the piezoelectric layers 105, and an inactive region B in which the internal electrode layers 107 of different polarities do not oppose each other via the piezoelectric layers 105. Accordingly, when the multi-layer piezoelectric element is operated, since only the active region A undergoes displacement and the inactive region B does not undergo displacement, stress may be concentrated in the border between the active region A and the inactive region B and become the start point of crack.\nAs shown in FIG. 5A, the multi-layer piezoelectric element 101 has the inactive layers 113 stacked on both ends thereof in the stacking direction. When the multi-layer piezoelectric element is operated, since the inactive region B does not undergo displacement, stress may be concentrated in the border between the region which undergoes displacement and the inactive layers 113 and become the start point of crack.\nSuch a crack as described above may grow from the border toward the side faces of the stack 109 (the inactive region B side), but also may grow toward the inside of the stack 109 (the active region A side). When an electric field is applied between the opposing internal electrode layers 107, the active layer A expands in the direction of the electric field by reverse piezoelectric effect and shrinks in the direction perpendicular to the direction of the electric field. When the piezoelectric layers 105 expand in the direction of the electric field, the element 101 expands in the stacking direction as a whole. In case the element 101 is housed in a casing or a frame which restricts the expansion, the element 101 receives a compressive force as a reaction.\nA crack which starts in the border and grows toward the active region A may run in the direction of thickness of the piezoelectric layers 105 while bending and branching in accordance to the state of stress. When a crack growing in the direction of thickness of the piezoelectric layers 105 is generated between the internal electrode layers 107 which oppose each other, there has been such a problem that short-circuiting occurs between the internal electrode layers 107, thus resulting in a decrease in the amount of displacement of the multi-layer piezoelectric element 101.\nIn recent years, it is a common practice to operate the multi-layer piezoelectric element continuously over a long period of time with a higher electric field applied, since it is required to achieve a large amount of displacement from a compact multi-layer piezoelectric element under a higher pressure. To meet these requirements, multi-layer piezoelectric elements having stress relieving layer provided inside thereof have been proposed (refer to, for example, DE10234787A1 and DE10307825A1). However, there have been demands for a multi-layer piezoelectric element which has higher durability under conditions of continuous operation over a long period of time with a high pressure."}
{"text": "Air conditioning units such as so-called \"window room air conditioners\" are commonly used for residential and similar applications and generally include closed refrigeration circuits having an evaporator and a condenser. The unit is normally divided by a partition into an evaporator section and a condenser section, the evaporator section communicating with the room air to be conditioned and the condenser section communicating with external air such as outdoor air. Refrigerant flows through the refrigerant circuit absorbing heat from the room air at the evaporator and discharging heat energy to the external air at the condenser. The conventional refrigeration circuit is completed by the addition of a compressor, an expansion device, and the appropriate connections between the components.\nSuch an air conditioning unit usually includes a basepan supporting all of the components and an outer housing surrounding the entire unit. The front of the evaporator, or indoor section, includes an indoor grille, which has openings therein for directing warm indoor air into the evaporator and discharge openings therein for directing air back into the room. The outdoor section of the housing includes a plurality of openings in the sides and top thereof, which serve as inlet openings for cooling air which flows into the outdoor section and outwardly therefrom after passing through the condenser coil, which is mounted vertically in the back of the outdoor section.\nIn addition to the components mentioned above, the outdoor section also typically includes an outdoor fan and fan orifice, as well as an electric motor, which typically also drives an indoor fan. The indoor section also typically includes the aforementioned indoor fan, an indoor fan orifice, a control box and a fan scroll structure for directing the air cooled by the evaporator back into the room to be cooled. Each of the aforementioned components requires means for attaching it to the basepan and/or other structure of the air conditioning unit. Numerous approaches are known for assembly of the components in a manner which will minimize the number of individual components required for assembly of the unit. The fewer components, and the fewer number of attachment means, results in lower material costs, less labor content and, accordingly, a less expensive unit."}
{"text": "Field of the Invention\nThe present invention relates to powder-based cosmetic compositions having improved performance characteristics. More specifically, the present invention is concerned with powder-based cosmetic compositions having a high affinity for keratinous surfaces which demonstrate improved transfer-resistance, long-wear, good coverage and oil-controlling properties. Moreover, the compositions feel silky and comfortable on the skin.\nDescription of the Prior Art\nWomen continually seek cosmetic products for application to skin, such as makeup foundation, pressed powder, blush and eyes shadow, that maintain their fresh look for extended periods of time and which exhibit transfer resistance. Powder-based cosmetics are well known. They are typically provided in solid form and comprise a major portion of powdered material, such as fillers and pigments, together with a binder, if the product is intended to be pressed into a container. Nevertheless, many such products tend to lack substantivity with the skin, and therefore are easily transferred to surfaces and must be reapplied frequently to main a fresh look. A major disadvantage with powder-based cosmetics which purport to provide long wear is that the extended wear is typically achieved by using solvents, waxes and/or polymers which adversely affect the feel of the product on the skin. Moreover, such ingredients can be irritating to skin. A further disadvantage of some known powder-based cosmetics is that they are composed substantially of talc. Talc is used as an absorbent, an opacifier, and to improve the feel of a product. Nevertheless, as an absorbent, talc not only soaks up oil, but, over time, also steals moisture from the product, causing it to lose its luxurious feel. Some products include silicon forming resins which can flake off. To address that problem, silicon-forming resins have been combined with plasticizers, but the resulting film may feel tacky on the skin.\nThere therefore remains a need for powder-based cosmetics which require fewer applications by consumers and have the attributes of a silky feel, long wear, transfer-resistance, resistance to streaking, longer lasting color and oil absorption, particularly under high humidity conditions."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a method and apparatus for wireless communication. More particularly, the present invention provides a method and apparatus for wireless communication utilizing impulse radio wherein an impulse radio communication system integrates with existing cellular/PCS wireless communication systems and/or current land line communication systems.\n2. Background of the Invention and Related Art\nThe demand for the ability to communicate voice and data in a mobile manner, not confined by the wires or fiber cables of non-wireless communication means, is ever increasing. Technology has improved in order to meet this demand, however, the current state of the art falls dramatically short in many important areas.\nAs anyone who has ever made a wireless call, whether it be from a xe2x80x9ccordless phonexe2x80x9d directly attached to land lines (land lines typically being twisted copper pairs and more recently fiber optic cables) or from a cellular/PCS phone, can attest to, the communication can be very unreliable. When using cordless phones, there is a severe limit on distance; further, quality can become a issue because of the multipath environment typical of most wireless communications. It is not uncommon to simply be xe2x80x9cdroppedxe2x80x9d from the communication or to hear xe2x80x9cnoisexe2x80x9d as you talk. Further, there are a limited number of channels that can be achieved because of frequency limitations. This inherently limits the number of users in a given area.\nA response to communication disruptions caused by problematic multipath environments has been to increase the transmit power. Previously the main problem with increasing transmit power was that in mobile communications, battery power is typically used. Therefore by increasing power, the battery life is diminished, thereby decreasing the mobility of the user (which is the primary purpose of wireless communications.) However, new studies have shown an even greater problem with existing cellular/PCS and wireless communication systems: it has been shown that the nonionizing electromagnetic radiation released by current wireless communications poses serious health threats to usersxe2x80x94to include cancer and brain tumors. Exacerbating this problem is the fact that power increases to overcome multipath problems cause an even greater health risk.\nThus, a great need exists in the wireless voice and data communication industry for a wireless communication means that not only overcomes many of the multipath problems associated with existing wireless communications, but does so at transmission power levels that diminish the health concerns now associated with existing wireless communications means.\nAs new technologies emerge, it is very important that they are capable of being integrated into existing technologies since vast amounts of resources are typically required to develop infrastructures. Such is especially the case for voice and data communications. Immense networks have been developed on both the wireless and wired fronts. Hundreds of thousands of miles of twisted copper pair and fiber optic cables have been laid throughout the world. Cellular and PCS base stations have been linked throughout the world and satellite communications are also integrated. Integration of some existing technologies has been accomplished; such as in U.S. Pat. No. 5,127,042, entitled, xe2x80x9cCellular cordless telephonexe2x80x9d, to Gillig etal. Gillig teaches integrating a cellular telephone system into an existing wireless telephone system. The integration of these two distinct systems accomplished its stated purpose, to wit, automatic transfer from the more expensive and less reliable, but more ubiquitous, cellular system to a more reliable and less expensive cordless system whenever practical. However, it does not accomplish a goal of diminishing the reliability and health concerns of both of the systems it is integrating.\nThus, not only is there a need to develop a wireless communication system that overcomes many of the multipath problems associated with existing wireless communications and at an extremely lower power level, but also to be able to integrate that system into existing communication technologies.\nIt is therefore an object of the present invention to provide a novel wireless voice and data communication means.\nIt is another object of the present invention to provide a wireless voice and data communication means which utilizes very low transmit powers.\nIt is a further object of the present invention to provide a wireless voice and data communication means which overcomes multipath problems associated with wireless transmissions.\nIt is another object of the present invention to provide a wireless voice and data communication means which utilizes very low transmit powers and overcomes multipath problems associated with wireless transmissions and integrates into existing voice and data communication infrastructures.\nThese and other objects are provided, according to the present invention, by a method and apparatus for wireless communication utilizing impulse radio, wherein an impulse radio communication system either stands alone or integrates with existing cellular/PCS wireless communication systems and/or current land line communication systems."}
{"text": "Emergency and medical personnel can be summoned in the United States to respond to an emergency situation by dialing “9-1-1” from any telephone. Although every effort is taken to minimize response time to the emergency situation, an appreciable amount of time may pass between the moment the initial call is made and when the emergency and medical personnel actually arrive at the scene of the emergency situation. The treatment window for certain health conditions (e.g., stroke) can be very small and it can be important to immediately administer medical aid while waiting for the emergency and medical personnel. However, the average person may not be knowledgeable about how to respond to various emergency situations.\nTherefore, what is needed is a system that can help the average person administer medical aid while waiting for the arrival of emergency and medical personnel."}
{"text": "Step-and-repeat projection microlithographic systems used in the microelectronics industry achieve important advantages, notably including extremely high quality. However, such systems also have major disadvantages. These include very low throughputs, and very high capital costs for the systems.\nIt would be a distinct benefit to the industry if a projection microlithographic system could be created that is characterized by high throughput and relatively low capital cost, yet still achieves a large part of the quality of step-and-repeat systems.\nOne significant reason why step-and-repeat systems are commonly used is that the substrates can shrink as the result of process steps. Shrinkage can be caused by such factors as the firing of the substrates, or the putting down of thin layers of materials on the substrates. Such layers may create stresses which pull on the substrates in such manner as to shrink them. Although both types of shrinkage are dimensionally small, they can create significant errors in the microelectronics components being fabricated. Step-and-repeat systems compensate for the shrinkage, and so reduce or eliminate these errors.\nPrior art related to projection microlithography has often involved Dyson projection or imaging systems. These 1-to-1 systems are well known to create important advantages. They have not, however, been previously so constructed or used as to achieve to a high degree (if at all), and/or in a high quality way, the following:\n(1) a large field size, such as 60 mm by 60 mm, with high quality, and with (for example) a 10-inch diameter primary mirror, so that throughput is massively increased; PA1 (2) large air gaps between the prisms (or mirrors) of the Dyson system, and the object and image planes thereof, in combination with a magnification adjustment element that is disposed in one of these air gaps and employed to compensate for such factors as the above-indicated shrinkage of the substrates; PA1 (3) an image plane that is parallel to the object plane, and an image that has the same orientation as the object, so that continuous scanning can be practically and effectively performed; PA1 (4) correction for the three near-UV spectral lines of mercury (365 nm, 405 nm, and 436 nm wavelengths) for the whole field, thus permitting much transfer of power from a mercury light source to the substrate, and correction for the alignment wavelength (577 nm) at the central, reduced, field; PA1 (5) a high numerical aperture (NA) for a 1-to-1 system, with consequent good resolution, and a large field size under these conditions; PA1 (6) magnification compensation or adjustment with very low distortion at all settings, and magnification compensation or adjustment with very little change in modulation transfer function at all settings; PA1 (7) Dyson system elements that are selected, constructed and located to perform functions including (a) correction of field aberrations, (b) correction of chromatic aberrations--with consequent good image quality over the large field and even with the large air gaps, and (c) reduction of astigmatism; PA1 (8) An adjustment mechanism for effecting very small adjustments in the positions of the prisms in order to aid in achieving precise registry or superposition of a point or small line in the projected image during one scan of a scanning projection microlithography system, in comparison to the image of the same point or line during the next scan."}
{"text": "This invention relates generally to marine propulsion devices, such as outboard motors, and, more particularly, to a drive shaft assembly for use in such a marine propulsion device.\nDrive shafts in conventional outboard motors are typically provided with one or more bearing races adjacent their lower ends and are exposed to corrosive water and exhaust gases at their upper ends. To provide improved lower end wear characteristics in conjunction with improved upper end corrosion resistance, prior drive shafts have been manufactured by inertially welding a lower drive shaft member of carbon steel to an upper drive shaft member of stainless steel. However, inertia welding is disadvantageous in that it requires costly secondary operations such as machining, heat treating and straightening.\nAttention is directed to the following U.S. Pat. No.:\nZackrisson 4,527,978 July 9, 1985 PA1 Federmann et al. 4,358,284 Nov. 9, 1982 PA1 Loyd et al. 4,272,971 June 16, 1981 PA1 Lakiza et al. 4,127,080 Nov. 28, 1978 PA1 Liaaen 4,112,866 Sept. 12, 1978 PA1 Campbell et al. 3,209,720 Oct. 5, 1965 PA1 Toelke 3,175,374 March 30, 1965 PA1 Shipley et al. 3,138,010 June 23, 1964 PA1 Willis, Jr., et al. 3,083,680 April 2, 1963"}
{"text": "Gas turbine engines may include a recuperator, which is a heat exchanger using hot exhaust gas from the engine to heat the compressed air exiting the compressor prior to circulation of the compressed air to the combustion chamber. Preheating the compressed air usually improves fuel efficiency of the engine. In addition, the recuperator reduces the heat of exhaust gas, which helps minimize the infrared signature of the aircraft.\nHowever, the high temperature differential within the recuperator typically induces thermal stresses which can significantly reduce its life. Seal failures may occur between the recuperator and the adjacent components, and such may affect the reliability of the engine. Some prior art recuperator structures are also relatively heavy and/or costly to manufacture, and as such generally unsuitable for aircraft use.\nAccordingly, improvements are desirable."}
{"text": "Electrified vehicles use electric machines for propulsion. Electric vehicles, hybrid electric vehicles (HEVs), and plug-in hybrid electric vehicles (PHEVs) are example types of electrified vehicles. Battery packs power the electric machines. Various factors impact decisions to initiate additional inspections, service, or replace components. Internal pressure within electrochemical battery cells of the battery packs is one such factor."}
{"text": "The present invention relates, in general, to a method of modulating physiological and pathological processes and, in particular, to a method of modulating cellular levels of oxidants and thereby processes in which such oxidants are a participant. The invention also relates to compounds and compositions suitable for use in such methods.\nOxidants are produced as part of the normal metabolism of all cells but also are an important component of the pathogenesis of many disease processes. Reactive oxygen species, for example, are critical elements of the pathogenesis of diseases of the lung, the cardiovascular system, the gastrointestinal system, the central nervous system and skeletal muscle. Oxygen free radicals also play a role in modulating the effects of nitric oxide (NO.). In this context, they contribute to the pathogenesis of vascular disorders, inflammatory diseases and the aging process.\nA critical balance of defensive enzymes against oxidants is required to maintain normal cell and organ function. Superoxide dismutases (SODs) are a family of metalloenzymes that catalyst the intra- and extracellular conversion of O2− into H2O2 plus O2, and represent the first line of defense against the detrimental effects of superoxide radicals. Mammals produce three distinct SODs. One is a dimeric copper- and zinc-containing enzyme (CuZn SOD) found in the cytosol of all cells. A second is a tetrameric manganese-containing SOD (Mn SOD) found within mitochondria, and the third is a tetrameric, glycosylated, copper- and zinc-containing enzyme (EC-SOD) found in the extracellular fluids and bound to the extracellular matrix. Several other important antioxidant enzymes are known to exist within cells, including catalase and glutathione peroxidase. While extracellular fluids and the extracellular matrix contain only small amounts of these enzymes, other extracellular antioxidants are also known to be present, including radical scavengers and inhibitors of lipid peroxidation, such as ascorbic acid, uric acid, and α-tocopherol (Halliwell et al., Arch. Biochem. Biophys. 280:1 (1990)).\nThe present invention relates generally to low molecular weight porphyrin compounds suitable for use in modulating intra- and extracellular processes in which superoxide radicals, or other oxidants such as hydrogen peroxide or peroxynitrite, are a participant. The compounds and methods of the invention find application in various physiologic and pathologic processes in which oxidative stress plays a role."}
{"text": "This invention relates to a mechanism for folding continuous-form sheeting, which is provided with perforations at predetermined intervals so that a part may be torn off easily. The sheeting is supplied in such a state that the sheets of the continuous-form sheeting have been stacked up by alternately folding them at the perforations, in such a way that the sheets of the continuous-form sheet may be stacked up by folding them at the perforations after the formation of images thereon in the same manner as before the formation of the images.\nHeretofore, continuous-form sheeting known as fan-folded sheeting has been employed in printers for printing, and particularly for outputs from computers. Such continuous-form sheeting is provided with perforations at predetermined intervals so that a part may be torn off easily and is supplied in such a state that sheets of the continuous-form sheeting have been stacked up by alternately folding them at the perforations.\nThe use of such continuous-form sheeting (for printing) affords numerous merits including: having output data checked in the order in which the data is output, facilitating proper arrangement of data as it is output continuously, making it possible to supply many sheets of forms without using a stacker, and simplifying sheet feed control, i.e., type position control, with accuracy by means of driving sprocket holes. Each of the sheets of the continuous-form sheeting after being used for printing, are stacked in line with their creases or perforations as before printing.\nThere have recently been developed imaging apparatus such as laser beam printers, the use of which is expanding steadily, for obtaining hard copies by utilizing electrophotography in a manner similar to electronic copying machines. In such machines, scanning the surface of a photoconductive drum is scanned and charged with a laser beam modulated according to data pertaining to images, such as graphic forms, characters and the like. The surface of the photoconductive drum is exposed to light, and the images are formed on recording sheeting.\nSuch a laser beam printer has been substantially patterned after the existing electronic copying machine, where a sheet cut in predetermined size, is used as the recording sheet. For a fixing device, a pair of fixing heat rollers are formed by disposing a backup roll in pressure-contact with a heat roller which is heated to a high temperature. In other words, the process generally adopted is a so-called fixation by means of a heat roller, in which recording sheeting carrying an unfixed toner image is passed between the rollers to fix the toner to the sheeting by fusion-bonding. In this case, there is an increased demand for the use of continuous-form sheeting.\nWhen continuous-form sheeting is employed in an imaging apparatus in which the fixing method by means of a heat roller is employed, the perforated creases where sheets are folded tend to become undone, because the continuous-form sheeting heated and pressed, i.e., heat-pressed in the fixing unit of the imaging apparatus. Thus, the continuous-form sheeting is not thereafter easily folded.\nThe arrangement stated above is therefore disadvantageous in that the continuous-form sheeting discharged from the imaging apparatus is left unfolded or irregularly bent, thus causing disorder in the vicinity of the sheet discharge port as the sheeting is irregularly bent.\nEven in the case of conventional printers such as wire dot printers and the like, difficulties arise in the folding and stacking of continuous-form sheeting after printing operations. For this reason, there has been contrived a means for guiding the sheet in a direction along which it is discharged in accordance with the bending direction thereof by providing a rocking guide arm, moving a stacker table back and forth in accordance with the bending direction of the sheeting or the like. However, these means proposed or implemented are invariably complicated in construction and not capable of effectively folding sheeting that has been heat-pressed during heat-roll fixation."}
{"text": "Counterfeiting of manufactured goods is a worldwide problem, with recent studies estimating that 8% of the world's total GDP is now generated by the manufacturing and sales of counterfeit products. Many classes of counterfeit goods create substantial risks to public health including counterfeit pharmaceutical drugs, auto parts, pesticides, and children's toys. In addition, counterfeit computer chips, aerospace parts, and identification documents present significant risks to national security.\nAuthentication alone is not enough to stop counterfeiting. Counterfeiters use a variety of strategies, including diverting unfinished products from factories that make authentic goods and then adding their own counterfeit brand identifiers such as labels and tags. Counterfeit items can enter the supply chain at any point, including at the original manufacturing facility, at the shipper, in distribution, or in retail stores. Unless the manufacturer or supplier can identify exactly where and when the item entered the supply chain, identifying and eliminating the counterfeit goods can be almost impossible.\nMany different approaches have been tried to uniquely identify and authenticate objects, including labeling and tagging strategies using serial numbers, bar codes, holographic labels, RFID tags, and hidden patterns using security inks or special fibers. All of these methods can be duplicated, and many add a substantial extra cost to the production of the goods sought to be protected. Physical labels and tags can also be easily lost, modified, or stolen."}
{"text": "The present invention relates to an ozone decomposing agent for decomposing ozone which is formed, for instance, in a corona charger for electrophotographic copying machines and laser printers, and in the apparatus for purifying water and air in which ozone is employed for such purification.\nOzone is formed by the action of sunlight, ultraviolet light or sparks of thunder on oxygen of the air. It is also formed when electrical apparatus for industrial use produces sparks in the air. For instance, a corona charger employed in copying machine ionizes the ambient air when it is in operation, and such ionization promotes the formation of ozone. Ozone is toxic to human beings and lower animals when breathing air containing more than 0.1 ppm of ozone for long periods of time. According to the safety standards for ozone to human beings proposed by Japan Association of Industrial Health, ACGIH (American Conference of Governmental Industrial Hygienists), and OSHA (Occupational Safety and Health Administration), the permissible maximum average concentration of ozone in the air is 0.1 ppm when breathing the air for 8 hours. Many apparatus for industrial use are produced by observing these standards. Ozone has a characteristic, pungent odor, and the odor is noticeable even at concentrations as low as 0.01 to 0.02 ppm. Some operators of an industrial apparatus which forms ozone at such concentrations may complain about the odor. When the concentration amounts to about 0.05 ppm, it has an unpleasant odor, and when the concentration exceeds 0.1 ppm, it is irritating to mucous membranes of the eyes and respiratory organs.\nFurther, ozone is a powerful oxidizing agent which oxidizes and deteriorates organic materials. Therefore, it is desirable that the concentration of ozone be as low as possible, not only to human beings, but also to industrial apparatus and devices.\nOzone is employed at concentrations as high as 500-2500 ppm for sterilization of water and for treatment of raw sewage by deodorization and decolorization thereof. For example, when water is sterilized by ozone, 1 to 3 g of ozone is blown into 1 m.sup.3 of water. Most of the ozone blown into water is decomposed in the water, but some of the remaining ozone is discharged from the water into the air. However, since the concentration of the thus discharged ozone in the air is usually as high as 1 ppm, it is necessary to decompose the discharged ozone before it spreads into the air for the safety to human beings and for the protection of environment.\nSince ozone is toxic to human beings when its concentration in the air is high, various methods have been proposed to decrease the concentration of ozone.\nFor example, filters made of an organic hydrocarbon such as activated carbon are employed for sewage purification.\nFilters comprising Hopcalite catalyst, such as metal oxides of manganese, copper, silver and cobalt, are also employed for decomposition of ozone. In addition to this, a method of using a paint containing a variety of organic materials has been proposed, although the decomposition efficiency is not high enough for use in practice.\nIn Japanese Laid-Open Patent Application No. 49-87334, there is proposed a method of decomposing ozone by using a rubber olefin polymer containing double-bond groups with which ozone may react, which method is for use in an electrophotographic copying machine.\nIn Japanese Laid-Open Patent Application No. 61-64315 there is proposed a method of removing ozone by using terpenoide.\nWhen any of filters containing activated carbon or Hopcalite catalyst is employed, if the density of the filter is high, the absorption of ozone thereto and the ozone decomposition efficiency can be increased. However, the flow rate of the air through the filter decreases and accordingly the pressure loss of the air increases. The result is that the temperature within the apparatus in which such filter is employed may increase and the life of the apparatus may be accordingly shortened. On the contrary, if the density of the filter is decreased, the absorption of ozone and the ozone decomposition efficiency are decreased. In any event, a suction device is necessary for causing the air containing ozone to pass through such filters, and in some cases, a cooling apparatus may become necessary for not raising the temperature within the apparatus. The result is that the entire price of the apparatus becomes high.\nWhen purifying water, a large quantity of activated carbon has to be used so that this method is expensive. Furthermore, the ozone decomposition power of activated carbon significantly decreases when activated carbon comes into contact with an organic solvent such as acetone.\nWhen a rubber olefin polymer is employed for decomposition of ozone, since the reaction of decomposition of ozone is of a solid-to-gas reaction type, the decomposition efficiency is not only low, but also it loses its decomposition power in a short time.\nIn addition to the above ozone decomposing agents, the following ozone decomposing agents are known: terpenoides capable of decomposing ozone, such as linalool, linalool ester and citral, dissolved in alcohol; natural essential oils containing, for instance, linalool; terpenoids; terpenoids absorbed in an absorbent such as porous silica powder. From these ozone decomposing agents, a terpenoide vaporizes and spreads into the air to react with ozone in the air, thereby decomposing ozone. In these ozone decomposing agents, however, the vaporization rate is not constant and it is difficult to continue the vaporization with a constant rate over a long period of time. More specifically, it has a tendency that it vaporizes very quickly in the beginning, and immediately the varporization drastically decreases. In order to continue the vaporization over a long period of time, a large amount of the ozone decomposing agents is required and therefore these ozone decomposing agents are expensive."}
{"text": "The invention relates to the field of cooling electronic devices and, in particular, to using circulating fluids to cool microprocessors, graphics processors, and other computer components.\nMicroprocessor dies typically used in personal computers are packaged in ceramic packages that have a lower surface provided with a large number of electrical contacts (e.g., pins) for connection to a socket mounted to a circuit board of a personal computer and an upper surface for thermal coupling to a heat sink. In the following description, a die and its package are referred to collectively as a microprocessor.\nElevation views of typical designs for heat sinks suggested by Intel Corporation for its Pentium® III microprocessor are shown in FIGS. 1A and 1B.\nIn FIG. 1A, a passive heat sink indicated generally by reference numeral 110 is shown. The passive heat sink 110 comprises a thermal plate 112 from the upper surface of which a number of fins, one of which is indicated by reference numeral 114, protrude perpendicularly. The passive heat sink 110 is shown in FIG. 1A installed upon a microprocessor generally indicated by reference numeral 118. The microprocessor 118 is comprised of a die 116 and a package 120. The die 116 protrudes from the upper surface of the package 120. The lower surface of the package 120 is plugged into a socket 122, which is in turn mounted on a circuit board (not shown). The passive heat sink 110 is installed by bringing the lower surface of the thermal plate 112 into contact with the exposed surface of the die 116. When installed and operated as recommended by the manufacturer, ambient airflow passes between the fins in the direction shown by an arrow 124 in FIG. 1A.\nIn FIG. 1B, an active heat sink, indicated generally by reference numeral 126, is shown. The active heat sink 126 comprises a thermal plate 128 from the upper surface of which a number of fins 130 protrude perpendicularly. A fan 132 is mounted above the fins 130. The active heat sink 126 is shown in FIG. 1B installed upon a microprocessor, generally indicated by reference numeral 136, which is comprised of a die 134 and a package 138. The die 134 protrudes from the upper surface of the package 138. The lower surface of the package 138 is plugged into a socket 140, which is in turn mounted on a circuit board (not shown). The active heat sink 126 is installed by bringing the lower surface of the thermal plate 128 into contact with the exposed surface of the die 134. When installed and operated as recommended by the manufacturer, ambient air is forced between the fins 130 in the direction shown by an arrow 142 in FIG. 1B.\nA difficulty with the cooling provided by the heat sinks shown in FIGS. 1A and 1B is that at best the temperature of the thermal plates 112, 128 can only approach the ambient air temperature. If the microprocessor 118, 136 is operated at a high enough frequency, the die 116, 134 can become so hot that it is difficult to maintain a safe operating temperature at the die 116, 134 using air cooling in the manner shown in FIGS. 1A and 1B.\nLiquid cooling, which is inherently more efficient due to the greater heat capacity of liquids, has been proposed for situations in which air cooling in the manner illustrated in FIGS. 1A and 1B is inadequate. In a typical liquid cooling system, such as that illustrated in FIG. 1C, a heat conductive block 144 having internal passages or a cavity (not shown) replaces the thermal plate 128 in FIG. 1B. The block 144 has an inlet and an outlet, one of which is visible and indicated by reference numeral 146 in FIG. 1C. Liquid is pumped into the block 144 through the inlet and passes out of the block 144 through the outlet to a radiator or chiller (not shown) located at some distance from the block 144. The block 144 is shown in FIG. 1C installed upon a microprocessor generally indicated by reference numeral 148, which is comprised of a die 150 and a package 152. The die 150 protrudes from the upper surface of the package 152. The lower surface of the package 152 is plugged into a socket 154, which is in turn mounted on a circuit board (not shown). The block 144 is installed by bringing its lower surface into contact with the exposed surface of the die 150.\nIn all liquid cooling systems known to the inventor, only a small portion of the lower surface of the block 144 comes into contact with the die 150. Since the die 150 protrudes above the upper surface of the package 152, a gap 156 remains between the upper surface of the package 152 and the block 144. If the gap 156 is not filled with insulation and sealed, convective and radiative heat transfer from the package 152 to the block 144 may occur. This will have no serious consequences so long as the block 144 is not cooled below the dew point of the air in the gap 156. If the liquid pumped through block 144 is only cooled by a radiator, then that liquid and consequently the block 144, can only approach the ambient air temperature. However, if a chiller is used to cool the liquid, then the temperature of the block 144 can decrease below the ambient air temperature, which may allow condensation to form on the package 152 or the block 144. Such condensation, if not removed, can cause electrical shorts, which may possibly destroy the microprocessor 148.\nCurrent solutions to the condensation problem referred to above include (1) controlling the chiller so that the temperature of the block 144 does not decrease below the dew point of the air in the gap 156 or (2) providing sufficient insulation and sealing material to prevent condensation from forming or to at least prevent any condensation that does form from reaching critical portions of the microprocessor 148 or surrounding circuit elements. Placing a lower limit on the temperature of the chiller limits the amount of heat that can effectively be removed from the microprocessor 148 without using bulky components. Further, the operating temperature of the microprocessor 148 can only approach the temperature of the block 144; operation at lower temperatures may be desirable in many circumstances. Alternatively, if insulation and sealing is used, trained technicians must do the installation properly if the installation is to be effective. If the insulation or seals fail, condensation can occur and cause catastrophic failure of the personal computer. A simpler, more reliable solution to the condensation problem is needed."}
{"text": "The filling of a large, moored floatation device and autonomous decoupling of the device from an umbilical fill line can require a system that can withstand high tensile loads resulting from drag and weight imparted by the mooring of a buoyant device (once filled). Design approaches in the prior art can use plug-and-socket style fill ports, with rotational actuating line releases. But these types of devices can snag or bind when subjected to loads greater than one hundred pounds force (100 lbf).\nIn view of the above, it can be an object of the present invention to provide a remotely actuated fill and release manifold that can withstand relatively high tensile forces once the buoyant fill bladder is filled. Another object of the present invention can be to provide a remotely actuated fill and release manifold that can allow for remote selectively release of a buoyant fill bladder from the device while the device under a relatively high tensile stress. Still another object of the present invention is to provide a remotely actuated fill and release manifold, which can allow for remote separation of the working components of the device without requiring the use of moving parts. Yet another object of the present invention to provide a remotely actuated fill and release manifold that that is easy to manufacture, that is inexpensive, and that is easy to deploy and use by remote operators in the field."}
{"text": "The present disclosure relates generally to the display of information, and more particularly to construction machine visualization.\nOperators of construction machines are required to operate the machines according to a desired site plan in order to modify the site to match the desired site plan. The operator is required to continuously monitor the current state of the site and operate the machine in a manner to change the current site to the desired site plan. Typically, an operator of a construction machine must constantly reference a desired site plan, typically on paper or on a display in the cab of the machine, to ensure that the modification of the current site is performed to change the site to the desired state."}
{"text": "This invention relates generally to apparatus and methods for animal semen collection, and more particularly to apparatus and methods using various techniques, including flow cytometry, to yield sperm populations that are enriched with sperm cells having one or more desired characteristics, such as viable populations of sperm cells sorted according to DNA characteristics for use by the animal production industry to preselect the sex of animal offspring.\nThe fertilization of animals by artificial insemination (AI) and embryo transplant following in vitro fertilization is an established practice. In the livestock production industry, the ability to influence the reproductive outcome toward offspring having one or more desired characteristics has obvious advantages. By way of example, there would be an economic benefit in the dairy industry to preselect offspring in favor of the female sex to ensure the production of dairy cows. Efforts have been made toward achieving this goal by using flow cytometry to sort X and Y sperm cells, as evidenced by the disclosures in U.S. Pat. No. 6,357,307 (Buchanan, et al.), U.S. Pat. No. 5,985,216 (Rens, et al.), and U.S. Pat. No. 5,135,759 (Johnson). However, none of these efforts has resulted in the introduction of a commercially successful high-throughput system capable of producing production volumes of relatively pure sexed sperm cells having a motility sufficient for effective fertilization.\nAccordingly, there is a current need in the animal production industry for a viable high-speed system for efficiently isolating sperm cells based on a specified DNA characteristic (or other characteristics) to produce quantities of such cells, which can be used on a commercial scale. Also needed is a sperm handling system that preserves the viability of such isolated sperm as it is processed by the isolating system and that allows for preservation of such isolated sperm until such time that it is ready for use. The present invention addresses these needs.\nThis invention also has application to improvements in the field of flow cytometry on a more general basis. Flow cytometry may broadly be defined as measuring characteristics of individual particles as they pass generally single file in a fluid stream through a measuring device which, typically, provides information for classifying the particles according to selected characteristics. Optionally, the particles may then be separated into populations using any number of techniques, including droplet sorting, droplet interference sorting, and fluid switching. Another option is to selectively destroy unwanted particles, for example by photo ablation.\nIn an optically-based flow cytometry system, optics are used to direct and focus a beam of light (e.g., visible light or UV light) on the stream containing the particles, and to collect emissions from the particles, including scattered light and/or fluorescence emissions from the particles. In one common optic system, for example, a beam of light (e.g., a laser beam) is focused on the stream and emissions are collected by a pair of collection units, one positioned forward of the laser for collecting scattered light emissions and another positioned orthogonally to both stream and the laser for collecting fluorescence emissions. Each collection unit includes a separate photodetector, which increases the cost of the system. Further, in traditional optic systems the photodetectors translate the collected emissions into electrical signals, which are analyzed using analog systems to classify the particles according to selected characteristics of the particles. Analog systems are relatively inexpensive, but only limited information can be derived from the signals.\nOthers have tried to develop technology that can be used to process sperm cells to obtain populations of sperm cells that are enriched with sperm that have a desired sex chromosome. However, the existing technology falls short of the inventive technologies described herein.\nFor example, Johnson et al. (U.S. Pat. No. 5,135,759) describe the separation of intact X and Y chromosome-bearing sperm populations according to DNA content using a flow cytometer/cell sorter into X and Y chromosome-bearing sperm enriched populations. As described, the sperm is combined with a DNA selective dye at a temperature of 30 to 39° C. for a period of 1 hour (39° C.) to 1.5 hours (30° C.). A flow cytometer is then used to measure the amount of fluorescent light emitted as the sperm passes through a laser beam that excites the dye. Because the X chromosome-bearing sperm contains more DNA than the Y chromosome-bearing sperm, with most species of mammal having about 3 to 5% difference, the X chromosome-bearing sperm emits more fluorescent light than the Y chromosome-bearing sperm. In order to account for the fact that the fluorescence measurement may vary depending on the rotational orientation of the sperm cells, two photo detectors are used. The first determines whether the sperm cells are properly oriented, while the second takes a measurement that is used to classify the sperm as having an X or Y chromosome. An oscillator is used to cause the stream containing the sperm to break into droplets downstream of the place where the sperm pass through the laser beam. Droplets containing single sperm of a predetermined fluorescent intensity are given a charge and electrostatically deflected into collection vessels. The collected, gender enriched sperm population, is then used for microinjection, in vitro fertilization, or artificial insemination.\nSeidel et al. (WO 02/43574) also describe separation of sperm into gender enriched populations of X and Y chromosome-bearing cells using flow cytometry. Seidel et al. describe staining the cells at a temperature between 30° C. and 40° C.\nUnited States Patent Application Pub. No. 2003/0157475 A1 (Schenk, Aug. 21, 2003) describes a method of cryopreserving sperm cells that have been sorted according to X or Y chromosome content. As noted therein, it is desirable to add a cryoprotectant to sperm cells before they are cryopreserved to protect the sperm cells during the cryopreservation process. For example, glycerol is one cryoprotectant that is commonly added to bovine sperm cells prior to cryopreservation. However, in order to obtain better protection from the cryoprotectant, it is desirable to wait for the cryoprotectant to equilibrate with the sperm cells before subjecting the sperm cells to temperatures below 0° C. During the equilibration period, the cryoprotectant penetrates the cell membrane to provide intra-cellular protection in addition to any extra-cellular protection provided by the cryoprotectant. Thus, the cryopreservation methods described in United States Patent Application Pub. No. 2003/0157475 A1 specify that an extender containing glycerol is added to the sperm cells after they have been cooled to about 5° C. Then the sperm cells and glycerol are allowed to equilibrate at 5° C. for anywhere between 1 and 18 hours before the sperm cells are subjected to lower temperatures. The disclosure recommends an equilibration period of between three and six hours in order to obtain the best results.\nUnfortunately, the time and expense involved in a 3 to 6 hour equilibration period will have a negative impact on profitability of a commercial sperm sorting process. Furthermore, in the context of a commercial sperm sorting process, it is believed that the health of the sperm is generally improved by reducing the time between collection of the sperm and cryopreservation (other factors being equal). From this standpoint as well, it would be desirable to have access to cryopreservation technology that does not require a long equilibration period to obtain the optimal benefits of a cryoprotectant. Moreover, the known cryopreservation technology is reported to have a detrimental impact on sperm motility, which is indicative of decreased sperm fertility. Thus, there is a need for cryopreservation techniques that preserves sperm health compared to conventional techniques."}
{"text": "1. Field of the Invention\nThis invention relates to fluid flow devices, particularly of that class of fluid flow nozzles with laminar discharge, and further to the field of illuminated fluid nozzles.\n2. Description of the Prior Art\nIn U.S. Pat. No. 5,160,086 granted on Nov. 3, 1992, a laminar flow fluid nozzle is provided for use in decorative water fountains, and also for industrial applications. Initial fluid flow is provided through a double-walled, bladder like fluid supply hose (32).\nThe inflated double hose isolates fluid flow from the pump and other vibrations which would otherwise be transmitted by means of a rigid hose. Furthmore, the double hose functions to absorb or accumulate small pressure variations known as \"pump noise\" which are usually present in the input stream.\nThe fluid next flows into a fluid chamber over a baffle (22). Inside pockets of air (18) cushion remaining pressure variations of the fluid stream.\nFluid next flows through a diffuser (20) which provides a purality of parallel fluid baths to dampen remaining major currents by lowering the fluid velocity, and thus the Reynolds number to obtain laminar fluid flow.\nA knife edged orifice (12) results in a narrow fluid stream (14). However, orifice (12) may be offset with respect to the center of enclosure (11).\nA light transmitting means (24) is aligned with orifice (12) by focusing means (28.\nA spring system (30) may be provided to absorb ambiant vibrations."}
{"text": "Individuals and organizations are rapidly accumulating large collections of digital image content, including still images, text, graphics, animated graphics, and full-motion video images. This content may be presented individually or combined in a wide variety of different forms, including documents, catalogs, presentations, still photographs, commercial videos, home movies, and meta data describing one or more associated digital content files. As these collections grow in number and diversity, individuals and organizations increasingly will require systems and methods for organizing and presenting the digital content in their collections. To meet this need, a variety of different systems and methods for organizing and presenting digital image content have been proposed.\nFor example, there are several manual digital image albuming systems that enable users to create digital photo albums manually. These systems typically provide tools for organizing a collection of images and laying out these images on one or more pages. Among the common types of tools for manually creating a digital photo album are tools for selecting a subset of images in the collection that will appear on a page of an album, a graphical user interface for manually rearranging the images on the page, and basic image editing tools for modifying various characteristics, such as size and orientation, of the images that will appear in the album. Users typically find the process of generating a digital photo album using fully manual digital image albuming systems to be tedious and time consuming.\nSome automated digital image albuming systems allow users to organize digital images into album pages in accordance with dates and times specified in the metadata associated with the images. These systems also typically allow users to annotate the images appearing in the digital photo album pages. Some automated digital image albuming systems provide various predefined layout templates that a user may select to create a digital photo album. In these systems, the user assigns images from the collection to various predefined image locations on a selected layout template, and the system automatically adjusts the size, placement, rotation, and framing of the images in accordance with parameters specified for the various predefined image locations on the selected template.\nWhat is needed are systems and method of arranging graphic objects on pages without predefined templates."}
{"text": "This application is related to our copending patent applications assigned to the assignee hereof:\nxe2x80x9cDIFFERENT WORD SIZE MULTIPROCESSOR EMULATIONxe2x80x9d by David A. Egolf, filed Sep. 8, 1998, with Ser. No. 09/149,261, and\nxe2x80x9cPROCESSOR EMULATION VIRTUAL MEMORY ADDRESS TRANSLATIONxe2x80x9d by David A. Egolf, filed of even date herewith, with Ser. No. 09/212,968.\nThe present invention generally relates to computer system emulation, and more specifically to instruction fetch optimization in emulation of a Target system utilizing a multiprocessor Host system.\nThe cost of designing a computer processor continues to increase. Some computer architectures thus ultimately become uneconomical to implement directly, despite these architectures having significant installed bases.\nOne solution to this problem is to simulate (or emulate) one computer architecture on another computer architecture. Herein, the simulating (or emulating) computer architecture will be termed the xe2x80x9cHostxe2x80x9d computer system, while the simulated (or emulated) computer architecture will be termed the xe2x80x9cTargetxe2x80x9d computer system. Emulators have been available almost since the advent of the first compilers.\nOne problem that arises in emulating one computer architecture on another involves address translation. In virtual memory computer systems, it is necessary to translate from operand and instruction virtual addresses to real memory addresses. In older computer systems, this translation was rudimentary. As computer technology progressed, this translation of virtual addresses to real addresses has continued to become ever more complicated.\nFull virtual to real address translation typically consumes significant computer cycles, whether done in hardware, or in an emulator. When emulating a complex computer architecture, such as the GCOS(copyright) 8 architecture provided to customers by the assignee herein, full virtual to real address translation may take upwards of a hundred instructions to accomplish. The problem is compounded by the necessity to translate at least two virtual addresses per instruction executed for most instructions executed or emulated: one to fetch the instruction to execute, and a second one for the instruction operand.\nHardware implementations have addressed this complexity and sped up the virtual to real address translations with a number of hardware assists that have been developed over time. One such hardware assist is the use of look-aside buffers, where recent operand and instruction addresses, along with their virtual to real address translations, are maintained. Instead of immediately doing virtual to real address translations, a search is made first of the look-aside buffer. If the virtual page address has been recently translated, the corresponding real page address for the operand or instruction can be taken from the look-aside buffer. This eliminates the repeated need for expensive virtual to real address translation. This ability to test or compare a number of addresses in parallel in hardware is well understood in the prior art, and is one of the fundamental bases for N-way set associative cache memories.\nSoftware emulations typically cannot take advantage of this look-aside buffer strategy since Host systems seldom supply general-purpose associative memory capabilities. Instead, they must resort to a costly search to implement the same strategy. A pragmatic result is that an implementation which performs the full address development is in many cases cheaper than implementing the look-aside buffer technique.\nWhen virtual to real address translation is done in software in an emulator, it is typically not possible on a Single Instruction/Single Data (SISD) processor to test or compare the virtual page address being translated against more than one previous virtual page address at a time. Thus, the previous virtual page addresses in the look-aside buffer would have to be tested sequentially. The result of this constraint is that it is often as expensive, or maybe even more expensive, in terms of computer cycles to utilize a look-aside buffer for caching virtual to real address translations as compared to performing full address translation for each address utilized.\nIt would thus be advantageous to be able to efficiently translate virtual to read addresses in an emulator. An efficient method of performing this address translation would significantly reduce the cost of emulating one (Target) computer system utilizing a second (Host) computer system."}
{"text": "In recent years, electronic devices have been used in each of which an electronic component that may be of various types is mounted on a circuit member of a heat dissipation base, examples of the component including semiconductor elements such as an insulating gate bipolar transistor (IGBT) element, a metal oxide semiconductor field effect transistor (MOSFET) element, a light emitting diode (LED) element, a freewheeling diode (FWD) element and a giant transistor (GTR) element, a sublimation-type thermal printer head element, and a thermal ink-jet printer head element.\nAs the heat dissipation base, on which a circuit member on which an electronic component is to be mounted is set up, there is used a heat dissipation base wherein a copper plate is joined, as a circuit member, to one of the main surfaces of a ceramic substrate and further a copper plate is joined, as a heat dissipation member good in heat dissipating property, to the other main surface. For joining the ceramic substrate to the circuit member and the heat dissipation member, a brazing material is used.\nFor this reason, the brazing material, which is used to join the circuit member and the heat dissipation member to the ceramic substrate, is required to have a high joining strength. The heat dissipation base, wherein the circuit member and the heat dissipation member are joined to the ceramic substrate through a joining layer made of the brazing material, is required to have an excellent heat dissipation property, and a high reliability that permits the base to be used for a long term.\nAs such a heat dissipation base, for example, Patent Literature 1 suggests a ceramic circuit board wherein metal circuit plates are joined to a ceramic substrate through a silver-copper based brazing material layer containing carbon powder and at least one active metal selected from Ti, Zr, Hf and Nb. Additionally, it is stated that when at least one element selected from In, Zn, Cd and Sn is further incorporated into the silver-copper based brazing material layer, the melting point of the brazing material is made low to lower the temperature for joining the metal circuit plates."}
{"text": "1. Field of the Invention\nThe present invention relates to an organic light emitting display (OLED), and more particularly, to an OLED that makes use of demultiplexers to decrease the number of output lines of data drivers and display an image having uniform luminance.\n2. Description of the Related Art\nIn recent years, a variety of flat panel displays (FPDs) with small weight and volume have been developed to overcome the drawbacks of a cathode ray tube (CRT). The FPDs can be categorized into a liquid crystal display (LCD), a field emission display (FED), a plasma display panel (PDP), an organic light emitting display (OLED), and the like.\nAn OLED includes an organic light emitting diode, which is a self-emissive display device that emits light through the recombination of electrons supplied from a cathode and holes supplied from an anode. The OLED has a high response speed of about 1 μs and consumes low power. The OLED includes a plurality of pixels, each of which supplies a driving current corresponding to a data signal to the organic light emitting diode using a driving thin film transistor (TFT) so that the organic light emitting diode emits light to display a predetermined image.\nFIG. 1 is a block diagram of a conventional OLED.\nReferring to FIG. 1, the conventional OLED includes a display panel 10, a scan driver 20, a data driver 30, and a timing controller 40.\nThe display panel 10 includes a plurality of pixels P11-Pnm, which are disposed in regions where a plurality of scan lines and emission control lines S1-Sn and E1-En intersect a plurality of data lines D1-Dm. Each of the pixels P11-Pnm receives a first power supply Vdd and a second power supply Vss from external power supplies and emits light corresponding to a received data signal to display an image. Also, the pixels P11-Pnm emit light for a time that is controlled according to a signal supplied through the emission control lines E1-En, respectively.\nThe scan driver 20 generates a scan signal in response to a scan control signal Sg received from the timing controller 40 and sequentially transmits the generated scan signal to the plurality of scan lines S1-Sn to select the pixels P11-Pnm. Also, the scan driver 20 generates an emission control signal in response to the scan control signal Sg and sequentially transmits the generated emission control signal to the plurality of emission control lines E1-En to control the emission.\nThe data driver 30 receives red (R), green (G), and blue (B) data from the timing controller 40, generates data signals in response to a data control signal Sd, and transmits the generated data signals to the plurality of data lines D1-Dm. In this case, the data driver 30 transmits data signals corresponding to one horizontal line for each horizontal period to the plurality of data lines D1-Dm.\nThe timing controller 40 generates data control signals Sd and scan control signals Sg corresponding to R, G, and B data supplied from an external graphic controller (not shown) and horizontal and vertical synchronous signals Hsync and Vsync. The data control signals Sd and the scan control signals Sg generated by the timing controller 40 are supplied to the data driver 30 and the scan driver 20, respectively.\nIn the conventional OLED having the above-described construction, the respective pixels P11-Pnm are located in the regions where the scan lines and the emission control lines S1-Sn and E1-En intersect the data lines D1-Dm. Here, the data driver 30 includes m output lines such that it transmits data signals to each of the m data lines D1-Dm. In other words, the data driver 30 should include the output lines in a number equal to the number of the data lines D1-Dm. Accordingly, the data driver 30 should include a plurality of data integrated circuits (ICs) to install the m output lines, thus elevating the cost of production. In particular, as the display panel 10 increases in resolution and size, the data driver 30 should include more data ICs. Thus, the cost of production further increases."}
{"text": "This invention relates to a macrolide composition. More particularly, this invention relates to a method of increasing the solubility of 3R, 4S, 5S, 6R, 7R, 9R, 11S, 12R, 13S, 14R-6-[(2S,3R,4S,6R)-4-dimethylamino-3-hydroxy-6-methyloxan-2-yl]oxy-14-ethyl-7,12,13-trihydroxy-4-[(2R,4R,5S,6S)-5-hydroxy-4-methoxy-4,6-dimethyloxan-2-yl]oxy-10-(2-methoxyethoxymethoxyimino)-3,5,7,9,11,13-hexamethyl-1-oxacyclotetradecan-2-one or roxithromycin. This invention further relates to a novel polymorph, (Form-III), of roxithromycin.\nRoxithromycin, a 14-membered-ring, macrolide antibiotic, is very similar in composition, chemical structure (semi-synthetic) and mechanism of action to erythromycin. Roxithromycin is currently available in two forms, namely anhydrous and monohydrate form.\nRoxithromycin exhibits activity against some sexually transmitted diseases, upper and lower respiratory tract infections, asthma, gum infections like gingivitis, and bacterial infections associated with stomach and intestinal ulcers. Roxithromycin is further regarded as the drug of choice for the treatment of opportunistic infections occurring in HIV/AIDS patients owing to its activity against Cryptosporidium spp., Mycobacterium avium complex, Pneumocystis carinii and Toxoplasma gondii. \nA disadvantage associated with roxithromycin is that it is a hydrophobic molecule because it has no free hydroxyl group and is thus poorly water-soluble and unstable in an acidic environment.\nA further disadvantage associated with roxithromycin is that its poor water-solubility and instability in an acidic environment result in a decrease in the absorption and bioavailability thereof.\nYet another disadvantage of roxithromycin is that said decreased absorption and bioavailability require relatively large quantities of roxithromycin to be administered in order to achieve a therapeutic effect.\nA disadvantage associated with the use of relative large quantities of roxithromycin is that there is a potential increase in the side-effects associated with this active ingredient.\nAn even further disadvantage associated with the use of relative large quantities of roxithromycin is that there is an increase in the production and manufacturing cost of roxithromycin dosage forms, thereby increasing the cost of treatment."}
{"text": "Many workers in the prior art have attempted to generate hydrogen and/or oxygen for use with an internal combustion engine. United States Patents Nos. 1,379,007 (Blumenburg); 1,520,772 (Ricardo); 3,648,668 (Pacheco); and 3,653,364 (Bogan) are all exemplary of efforts made in this direction. The United States Patent to Boisen 1,380,183 of May 31, 1921, is a teaching of an early effort to generate oxygen and hydrogen and use both gases to fuel an engine. Boisen also uses separate tanks in which to create a back pressure to force the electrolyte away from both the cathode and the anode to proceed for a stop and start of the decomposing process automatically. Boisen is a three-compartment unit utilizing the hydrogen and oxygen. Many workers in the art, however, believe the presence of storage tanks such as 11 and 12 for the hydrogen and oxygen of Boisen and the subsequent mixing of these fuels, presents a highly dangerous environment. The United States patent to Jack H. Rupe No. 3,906,913 (1975) also teaches the mixing of hydrogen, air and a liquid hydrocarbon fuel. The patent is particularly useful in its description of the prior art and its disclosure of the degree of work some workers in the prior art did to secure proper levels of hydrogen. Another patent of interest is the United States patent to Marc S. Newkirk, et al. 3,710,770, of June 16, 1973. Newkirk discloses a means of providing hydrogen from a cryogenic supply.\nA principal objective of this invention is to provide a rugged producer of hydrogen gas for an internal combustion machine having a pair of compartments with anodes and cathodes therein. The invention includes a means for causing a back pressure against the level of electrolyte in said cathode compartment to lower the level of electrolyte therein and thereby reduce the surface area available for the production of hydrogen. In this manner, the necessity for separate storage tanks or vessels is eliminated.\nAnother principal objective of this invention is to provide a rugged, uncomplex, two-compartment reactor for developing adjustable amounts of hydrogen. The reactor eliminates the necessity of storing any dangerous amounts of hydrogen in the system without a requirement for complicated valving and measuring systems. Except for the small amount of hydrogen necessary to build a back-pressure for adjustability, the hydrogen is used as it is produced.\nAn important object of the present invention is to provide a hydrogen reactor for economically providing on demand the trace amounts of hydrogen necessary to operate an internal combustion engine with a mixture of said hydrogen, air and a hydrocarbon-type fuel.\nA still further important objective of the invention is to provide a hydrogen generator having a triangular-shaped electrode the surface, so that the incremental amount of active surface area exposed or covered by each successive incremental fall or rise in the electrolyte is smaller or larger, respectively, than the incremental amount of active surface area exposed or covered by the previous incremental fall or rise in accordance with the demands from the engine. The generator is also equipped with a separator that divides the generator into an electrode chamber and an anode chamber. The divider is of a composition that substantially limits the flow of electrolyte between chambers but which transmits electrons therebetween.\nAnother important objective of this invention is to provide an internal combustion engine which, because of superior burning caused by the hydrogen, permits the ignition timing to be advanced from that in conventional ignition systems.\nAnother principal objective of the invention is to improve the quality of the atmosphere, especially in areas where large numbers of vehicles are used by releasing substantial amounts of oxygen into the atmosphere.\nA still further objective of the invention is to provide a fuel generator and mixture system which is readily adapted for use with all engines and of particular usefulness with burners of gas turbine engines and pre-compression chambers of two-stroke motors.\nThe present invention relates to a process for operating an internal combustion engine with a mixture of hydrogen, air and a hydrocarbon fuel, comprising the step of producing and introducing trace amounts of hydrogen into a carburetor together with said hydrocarbon fuel and air in adjustable amounts dependent on the hydrogen requirements of the engine.\nIt is well known in the art to mix hydrogen with a mixture of gasoline vapor and air produced in the carburetors of internal combustion engines to enhance the efficiency of such engines. It is also known that a better and more complete oxidation of the fuel in the combustion chamber of the internal combustion engine can be obtained, when said engine is operated with a lean air-fuel mixture. This has, however, the drawback that for conventional hydrocarbon type fuels the combustion of weak mixtures, i.e., of mixtures having a proportion of air significantly larger than that corresponding to the stoichiometric fuel/air ratio, may lead to misfire, uncontrolled combustion and possibly breakdown of the engine. On the other hand, a very weak mixture of hydrogen and of air is easily ignited. Excellent inflammability is characteristic for a hydrocarbon-fuel/trace amounts of hydrogen/air mixture. When such a mixture is ignited by means of the spark plug, the hydrogen which distributes itself throughout the mixture burns first and contributes to an optimum combination of the remaining vaporized fuel.\nBy using the aforementioned principles, it is possible to enlarge the ignition range of the engine and to burn mixtures having a high proportion of air. The temperature of the combustion is lower due to the additional air. Heat losses are accordingly reduced. Such a nearly perfect combustion in combination with high compression ratios improves the economy of the engine, reduces the production of harmful exhaust gases and delays the accumulation of soot in the combustion chamber of the engine. Taking advantage of these known principles has been difficult because no appropriate device has been available for producing and for carrying hydrogen for mobile internal combustion engines. Taking along heavy pressure bottles for hydrogen gas is not convenient, because such bottles must be filled up or exchanged quite frequently. In addition, such pressure bottles present the danger of explosion. Another possibility, the stocking of liquefied hydrogen at low temperatures is troublesome, expensive and dangerous, especially for mobile use.\nIn the process for operating an internal combustion engine, according to the present invention, hydrogen is produced by electrolytic decomposition of water in a hydrogen reactor, electric current is supplied to said hydrogen generator by means of an alternator, and the hydrogen generator reactor automatically adjusts its hydrogen production to the requirements of the engine. According to the present invention, the automatic adjustment of the hydrogen production rate of said hydrogen generator to the needs of the internal combustion engine is accomplished by causing hydrogen collected between the carburetor and the reactor to exert a pressure on the level of the electrolyte in such a manner that the electrolyte level is lifted or lowered in accordance with the hydrogen pressure of the hydrogen reactor, whereby the surface area of the cathode plates in the electrolyte is adjusted. In this manner, the level of the electrolyte is caused to drop in the cathode compartment when the hydrogen pressure in the hydrogen generator reaches a predetermined magnitude. The electric current of the hydrogen generator is thus automatically lowered and the production of hydrogen in the hydrogen generator is reduced. The hydrogen reactor and its plates are triangularly shaped to increase the sensitivity of the reactor to fuel demands.\nThe present invention comprises a self-controlling hydrogen reactor operating according to an electrolytic type process, and a DC generator for supplying current to said hydrogen generator. The current generator for supplying current to said hydrogen generator is driven by the engine of the mobile unit."}
{"text": "Ultrasound imaging is a well established tool for medical diagnosis and nondestructive testing and inspection of materials. Ultrasound imaging systems use transducers to create short high frequency acoustic pulses. The acoustic pulses propagate into the object under investigation. At locations where there is a change in acoustic properties, such as the boundary between two different tissue layers, part of the ultrasound energy is reflected. The reflected echos are detected by the transducer and processed to produce a two dimensional image of the underlying structures. The ability of the imaging system to detect small or subtle structures is primarily determined by the ability of the transducer to focus the ultrasound energy. Focusing the ultrasound beam can be achieved by shaping the transducer, by using an acoustic lens, or by using an array of transducers. Most modern ultrasound imaging systems employ transducer arrays.\nThe ultrasound energy produced by an array is focused by introducing time delays to the signals delivered to (or received from) individual array elements so that signals transmitted to (or received from) the desired region in space constructively interfere while signals outside this region destructively interfere. How well an array achieves this constructive interference is determined by the radiation pattern of the array. The radiation pattern can be visualized as a plot of the amplitude of the signal transmitted or received by the array as a function of position in space. An example of a radiation pattern is shown in FIG. 1, where the radiation pattern has been plotted as a contour plot using a logarithmic scale (dB) normalized to the peak pressure. The peak in the radiation pattern corresponds to the desired region in space over which the energy will be transmitted and from which the energy will be received. The width of the main peak in the radiation pattern is inversely proportional to the width of the array and determines the resolution of the imaging system. The non-zero amplitude of the radiation pattern away from the peak is caused by imperfect destructive interference and results in transmission and reception of unwanted energy. Transducer arrays are designed with the goal of obtaining a peak in the radiation pattern that is as narrow as possible while minimizing the amplitude of energy in the radiation pattern away from the main peak.\nUltrasound transducer arrays are fabricated by cutting a series of narrow grooves or kerfs through a bulk transducer substrate such as lead zirconate titanate (PZT). The grooves are used to mechanically and electrically isolate the array elements. To provide mechanical support to the narrow array elements, the grooves are often filled with a soft polymer material. The resulting array is essentially a composite structure consisting of alternating layers of PZT ceramic and polymer. Transducer arrays have also been fabricated by forming an electrode pattern on the surface of a ceramic-polymer composite fabricated using other methods, or on a polymer transducer substrate such as poly(vinylidene fluoride) (PVDF).\nThe grooves or kerfs of most ceramic arrays are machined using a thin diamond wheel. Other machining techniques such as laser machining or ultrasonic machining have also been used to separate the array elements. Although these techniques work well for arrays designed to operate below 10 MHz, they are not suitable for machining the extemely small and closely spaced elements required for high frequency imaging. They are also difficult to use with new single crystal relaxor ferroelectric materials such as lead zirconate niobate-lead titanate (PZN-PT) and lead magnesium niobate-lead titanate (PMN-PT) which are brittle and prone to chipping and cracking."}
{"text": "More or less unauthorized copying of media content that is copyright protected has been an important issue for a long period of time. For example, audio records have been copied onto cassette tapes and television programs and films have been copied onto video tape. Unauthorized copies of such types of analog recording media have all been marred by at least a slight decrease of quality as compared to the original recordings.\nHowever, the issue has become even more important as the media industry has entered “the digital era”. Music and films are now being distributed in digitally encoded form and copies of these are usually indistinguishable from the original, creating an even bigger market for “pirate” of media content.\nTo illustrate the situation one can consider the recordable compact disc, i.e. the CD-R. These are usually intended for recording audio content. However, any media content may be recorded on such discs. Moreover, while the following discussion will be exemplified by CD-R, there exist other digital recording media capable of recording media in the form of, e.g., films that are equally relevant for the discussion.\nEver since its market introduction, CD-R has been an overwhelming success story. One of the main factors contributing to its rapid acceptance is that consumers now have the possibility to make digital quality copies of the most popular audio format—the Compact Disc—that will play in any CD player. Whereas CE style CD recorders, i.e. dedicated digital audio copying devices, have to implement the Serial Copy Management System (SCMS), which limits duplication of copyright asserted discs to first generation copies, there are no built-in limitations for PC based CD recorders. The latter has become a huge problem from the perspective of the content owners, which are facing the horror scenario that immediately after release of, e.g., a new music album, large numbers of CD-R copies are traded or offered for sale on school yards and similar environments. Consequently, and especially in combination with the mp3 phenomenon on the Internet, content owners consider CD-R a major factor in the stagnating growth of CD sales.\nTechnical measures that prevent misuse of recorders, such as making large numbers of copies are therefore highly desirable. At the same time the convenience of the devices should not be compromised, e.g. if copying is allowed by the owner of the content, then one should be able to do so at high speed.\nA prior art solution within the technical field of copy protection of digital data is presented in the International Patent Application published with the number WO 00/04688. There, a method and a system for protecting digitally watermarked media content includes generating digital tickets associated with the media content. The ticket holds information regarding a receiver and the copy protection status of the content as well as the digital watermark. The content of the ticket is processed by the receiver in order to determine whether or not the receiver is allowed to copy the content."}
{"text": "1. Field of the Invention\nThe present invention is directed to a computer-implemented product for locating and connecting to a particular desired object or target resource from among plural resources resident at distributed locations on a network.\n2. Description of the Related Art\nThe worldwide network of computers known as the Internet evolved from military and educational networks developed in the late 1960's. Public interest in the Internet has increased of late due to the development of the World Wide Web (hereinafter, the Web), a subset of the Internet that includes all connected servers offering access to hypertext transfer protocol (HTTP) space. To navigate the Web, browsers have been developed that give a user the ability to download files from Web pages, data files on server electronic systems, written in HyperText Mark-Up Language (HTML).\nWeb pages may be located on the Web by means of their electronic addresses, known as Uniform Resource Locators (URLs).\nA URL uniquely identifies the location of a resource (web page) within the Web. Each URL consists of a string of characters defining the type of protocol needed to access the resource (e.g., HTTP), a network domain identifier, identification of the particular computer on which the resource is located, and directory path information within the computer's file structure. The domain name is assigned by Network Solutions Registration Services after completion of a registration process.\nWhile the amount of information available on the Web is enormous, and therefore potentially of great value, the sheer size of the Web makes the search for information, and particular web sites or pages, a daunting task. Search engines have been developed to assist persons using the Web in searching for web pages that may contain useful information.\nSearch engines fall into two major categories. In search engines falling into the first category, a service provider compiles a directory of Web sites that the provider's editors believe would be of interest to users of the service. The Yahoo site is the best known example of such a provider. Products in this category are not, strictly speaking, search engines, but directories, and will be referred to hereinafter as “editor-controlled directories”. In an editor-controlled directory, the developer of the directory (the “editor”) determines, based upon what it believes users want, what search terms map to what web pages.\nThe other major category, exemplified by Altavista, Lycos, and Hotbot, uses search programs, called “web crawlers”, “web spiders”, or “robots”, to actively search the Web for pages to be indexed, which are then retrieved and scanned to build indexes. Most commonly, this is done by processing the full text of the page and extracting words, phrases, and related descriptors (word adjacencies, frequencies, etc.). This is often supplemented by examining descriptive information about the Web document contained in a tag or tags in the header of a page. Such tags are known as “metatags” and the descriptive information contained therein as “metadata”. These products will be referred to hereinafter as “author-controlled search engines,” since the authors of the Web documents themselves control, to some extent, whether or not a search will find their document, based upon the metadata that the author includes in the document.\nEach type of product has its disadvantages. Author-controlled search engines tend to produce search results of enormous size. However, they have not been reliable in reducing the large body of information to a manageable set of relevant results. Further, web site authors often attempt to skew their site's position in the search results of author-controlled search engines by loading their web site metatags with multiple occurrences of certain words commonly used in searches.\nEditor-controlled directories are more selective in this regard. However, because conventional editor-controlled directories do not actively search the web for matches to particular search terms, they may miss highly relevant web sites that were not deemed by the editors to be worthy of inclusion in the directory. Also, it is possible for the editor to “play favorites” among the multitude of Web documents by mapping certain Web documents to more search terms than others.\nRecently, search engines such as DirectHit (www.directhit.com) have introduced feedback and learning techniques to increase the relevancy of search results. DirectHit purports to use feedback to iteratively modify search result rankings based on which search result links are actually accessed by users. Another factor purportedly used in the DirectHit service in weighting the results is the amount of time the user spends at the linked site. The theory behind such techniques is that, in general, the more people that link on a search result, and the longer the amount of time they spend there, the greater the likelihood that users have found this particular site relevant to the entered search terms. Accordingly, such popular sites are weighted and appear higher in subsequent result lists for the same search terms. The Lycos search engine (www.lycos.com) also uses feedback, but only at the time of crawling, not in ranking of results. In the Lycos search engine, as described in U.S. Pat. No. 5,748,954, priority of crawling is set based upon how many times a listed web site is linked to from other web sites. This idea of using information on links to a page was later exploited by the Clever system developed in research by IBM, and the Google system (www.google.com), which do use such information to rank possible hits for a search query.\nEven leaving aside the drawbacks discussed above, search engines of both categories are most useful when a user desires a list of relevant web sites for particular search terms. Often, users wish to locate a particular web site but do not know the exact URL of the desired web site. Conventional search engines are not the most efficient tools for doing this.\nMoreover, naming and locating particular sites on the Web is currently subject to serious problems. For example, appropriate names, including existing company names or trademarks, may not be available, because someone registered them first. Names may be awkward and not obvious, because of length, form/coding difficulties or variant forms, and names may not justify a separate domain name registration for reasons of cost and convenience, such as movie titles or individual products.\nThis problem results from a mismatch between the present network addressing scheme based on Uniform Resource Locators (URLs), which meet the technical needs of the Internet software, and the needs of human users and site sponsors for simple, user-friendly mnemonic and branded names. This problem is largely hidden in cases where a user finds a site by clicking a pre-coded link (such as after using a search engine), or by using a saved bookmark. However, the problem does seriously affect users wishing to find a site directly, or to tell another person how to find it. To do this, the person must know and type the URL into his Internet browser, typically of the form sitename.com or www.sitename.com. Site sponsors are also seriously hampered by this difficulty in publicizing their sites.\nFurther, the current method of naming and locating Web sites has serious, widely known problems. Web site locator “domain” names are often not simple or easily remembered or guessed, and often do not correspond to company, trademark, brand or other well-known names.\nAs a result of the foregoing, site URLs (or domain names) are not intuitively obvious in most cases, and incorrect access attempts waste time and produce cryptic error messages that provide no clue as to what the correct URL might be. A significant percentage of searches are for specific, well-known sites. These could be found much more quickly by a special-purpose locator engine. The current mode of interacting with search engines is also cumbersome-for this purpose, a much simplified mode of direct entry is practical.\nOne attempt to provide the ability to map a signifier, or alias, to a specific URL utilizes registration of key words, or aliases, which when entered at a specified search engine, will associate the entered key word with the URL of the registered site. One such commercial implementation of this technique is known as NetWord (www.netword.com). However, the NetWord aliases are assigned on a registration basis, that is, owners of web sites pay NetWord a registration fee to be mapped to by a particular key word. As a result, the URL returned by NetWord may have little or no relation to what a user actually would be looking for. Another key word system, RealNames (www.realnames.com), similarly allows web site owners to register, for a fee, one or more “RealNames” that can be typed into browser incorporating RealNames' software, in lieu of a URL. Since RealNames also is registration based, there is no guarantee that the URL to which is user is directed will be the one he intended.\nFurther, in existing preference learning and rating mechanisms, such as collaborative filtering (CF) and relevance feedback (RF), the objective is to evaluate and rank the appeal of the best n out of m sites or pages or documents, where none of the n options are necessarily known to the user in advance, and no specific one is presumed to be intended. It is a matter of interest in any suitable hit, not intent for a specific target. Results may be evaluated in terms of precision (whether “poor” matches are included) and recall (whether “good” matches are not included).\nA search for “IBM” may be for the IBM Web site, but it could just as likely be for articles about IBM as a company, or articles with information on IBM-compatible PCs, etc. Typical searches are for information about the search term, and can be satisfied by any number of “relevant” items, any or all of which may be previously unknown to the searcher. In this sense there is no specific target object (page, document, record, etc.), only some open ended set of objects which may be useful with regard to the search term. The discovery search term does not signify a single intended object, but specifies a term (which is an attribute associated with one or more objects) presumed to lead to any number of relevant items. Expert searchers may use searches that specify the subject indirectly, to avoid spurious hits that happen to contain a more direct term. For example, searching for information about the book Gone With The Wind may be better done by searching for Margaret Mitchell, because the title will return too many irrelevant hits that are not about the book itself (but may be desired for some other task).\nIn other words, the general case of discovery searching that typical search engines are tuned to serve is one where a search is desired to return some number, n, of objects, all of which are relevant. A key performance metric, recall, is the completeness of the set of results returned. The case of a signifier for an object, is the special case of n=1. Only one specific item is sought. Items that are not intended are not desired—their relevance is zero, no matter how good or interesting they may be in another context. The top DirectHit for “Clinton” was a Monica Lewinsky page. That is probably not because people searching for Clinton actually intended to get that page, but because of serendipity and temptation—which is a distraction, if what we want is to find the White House Web site.\nIn addition,                CF obtains feedback from a group of users in order to serve each given user on an overall, non-contingent basis—without regard to the either the intent of the user at a specific time, or to being requested in a specific context.        RF is used by a single user to provide feedback on their intent at a given time, but still with no presumed intent of a single target.        \nMore broadly, searching techniques are generally not optimized based on using a descriptor which is also an identifier—they provide more generally for the descriptor to specify the nature of the content of the target, not its name. There are options in advanced search techniques which allow specification that the descriptor is actually an identifier, such as for searching by title. Such options may be used to constrain the search when a specific target happens to be intended, but no special provision is made to apply feedback to exploit that particular relationship or its singularity.\nMoreover, none of the currently available key word systems utilize heuristic techniques actually to determine the site intended by the user. Instead, the current systems teach away from such an approach by their use of registration, rather than user intention, to assign key words to map to web pages. Thus, the current techniques are not directed to solving the problem of finding the one, correct site for a particular signifier.\nThus, the need exists for a system that would enable a user to find a desired Web document by simply entering an intuitive key word or alias and that would perform a one to one mapping of the alias with the URL actually desired by the user, and which would use heuristic techniques to assist in providing the correct mapping, and improving system accuracy over time."}
{"text": "Large wall displays can be prohibitively expensive as the cost to manufacture display panels rises exponentially with display area. This exponential rise in cost arises from the increased complexity of large monolithic displays, the decrease in yields associated with large displays (a greater number of components must be defect free for large displays), and increased shipping, delivery, and setup costs. The exponential rise in cost also includes the energy requirements associated with large wall displays. As a result, typical large wall displays are impractical for personal or home usage."}
{"text": "1. Field of the Invention\nThe present invention relates to a power transmission apparatus using gears. More particularly, it relates to a power transmission apparatus having a first shaft and a second shaft each supported by an aluminum housing, and an iron gear and a synthetic resin gear which are mounted on to said two shafts respectively and are engaged with each other.\n2. Description of the Prior Art\nPreviously known materials used for a gear, a power transmission element, include those which comprise a synthetic resin as the main component (see, for example, Japanese Utility Model Application Kokai (Laid-Open) No. 70,702/85). Synthetic resin gears of this kind are used, for example, for camshaft gears of engines. An example of such uses will be described below with reference to the accompanying figure illustrating a working example of that utility model.\nIn the FIGURE, 1 indicates the crankshaft of an engine (a first shaft), 2 indicates the camshaft (a second shaft), 4a indicates a crankcase (housing) and 4b indicates the crankcase cover (housing). The crankshaft 1 and the camshaft 2 are arranged in parallel with each other and are supported by the crankcase 4a and the crankcase cover 4b. The crankshaft gear 5 made of iron press-fitted to the crankshaft 1 and the camshaft gear 7 made of synthetic resin fixed to the camshaft 2 with rivets 6 engage with each other and transmit the torque of the crankshaft 1 to the camshaft 2. In the prior art, as the material for the camshaft gear 7, synthetic resins containing with glass fibers have been used to improve the mechanical properties of the camshaft gear 7.\nSince the camshaft gear 7 is made of a synthetic resin, it is abraded by meshing with the crankshaft gear 5 made of iron even though the resin contains glass fibers, causing the glass fibers to fall off into lubricating oil. Since the glass fibers which fall off are minute, they get into, for example, such sliding parts as the bearing parts 8 and 9 of the crankshaft 1 and the camshaft 2. Since the glass fiber is very hard, the above-mentioned phenomenon brings about the disadvantage of abrasion of the sliding parts. Further, there is another disadvantage in that the glass fibers come up to the surface of the gear teeth on the side of the camshaft gear 7 made of resin, whereby the abrasion of the opposing crankshaft gear 5 of iron is increased.\nAs a means for solving such problems, it is possible to secure the necessary mechanical properties of the gear, without including glass fibers in the synthetic resin, by increasing the face width, changing to a resin with more enhanced heat resistance, or adding a substance capable of enhancing sliding such as molybdenum. However, the elimination of glass fibers increases the coefficient of linear expansion of the resulting material and markedly increases the change of backlash with temperature. Thus, the coefficient of linear expansion is 2.3.times.10.sup.-5 /deg C for aluminum, 1.2.times.10.sup.-5 /deg C for iron and 9.times.10.sup.-5 /deg C for the synthetic resin. Accordingly, with respect to the distance D between the crankshaft 1 and the crankshaft 2, the sum of the magnitude of expansion (or contraction) of the two gears 5 and 7 is larger than the magnitude of expansion (or contraction) of the crankcase 4, so that the backlash between the two gears 5 and 7 will differ greatly from the set value with changes in temperature.\nIf the change of backlash with temperature is large, when, for example, the backlash is so set as to be optimum at high temperatures, the backlash will become excessively large at low temperatures, causing emission of noise. On the contrary, when the backlash is so set as to be optimal at low temperatures, the backlash will become small at high temperatures, causing emission of noise and severe abrasion.\nOn the other hand, phenolic resins and the like containing cloth or wood chips have little of the above-described problems with respect to both contamination and backlash. However, generally they require the process of gear cutting, which markedly lowers the production efficiency."}
{"text": "1. Field of the Disclosure\nThe present invention relates generally semiconductor processing. More specifically, examples of the present invention are related to semiconductor processing of stacked integrated circuit systems.\n2. Background\nAs integrated circuit technologies continue to advance, there are continuing efforts to increase performance and density, improve form factor, and reduce costs. The implementation of stacked three dimensional integrated circuits have been one approach that designers sometimes use to realize these benefits. The advances in wafer bonding with very precise alignments make it possible to fabricate stacked chips on wafer-level. The possible applications could include logic chip bonding to memory, image sensors, among others. This offers the advantage of smaller form factor, higher performance, and lower cost.\nA key challenge when implementing stacked three dimensional complementary metal-oxide semiconductor (“CMOS”) image sensors, which continue to get smaller and faster, relates to switching noise between the stacked three dimensional integrated circuits. Switching noise can be of particular concern to image sensors packaged in through-silicon vias (“TSV”) technology. With such packages, a number of traces or signal lines are laid out on the bottom side of the package. These traces often connect vias on the outer perimeter to solder balls (pins) in the inner region. During sensor operation, if a pin switches rapidly between high low states, and its corresponding trace runs underneath a sensitive part of the image sensor (e.g., pixel array), then switching noise may be coupled into the image sensor circuitry. This coupled noise may degrade the quality or increase noise in the output image data. The noise contributed from a pin depends on the location of the trace, the run length below the image sensor, the frequency of the switching, and the current in the trace. However, the noise emanating from these traces can affect a portion of the image sensor and even potentially the whole image sensor.\nCorresponding reference characters indicate corresponding components throughout the several views of the drawings. Skilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present invention. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments of the present invention."}
{"text": "The invention relates to the general field of magnetic disks with particular reference to GMR read heads and use of longitudinal bias leads therewith.\nThe present invention is concerned with the manufacture of the read element in a magnetic disk system. This is a thin slice of material, located between two magnetic shields, whose electrical resistivity changes on exposure to a magnetic field. Magneto-resistance can be significantly increased by means of a structure known as a spin valve (SV). The resulting increase (known as Giant magneto-resistance or GMR) derives from the fact that electrons in a magnetized solid are subject to significantly less scattering by the lattice when their own magnetization vectors (due to spin) are parallel (as opposed to anti-parallel) to the direction of magnetization of the solid as a whole.\nThe key elements of a spin valve structure are two magnetic layers separated by a non-magnetic layer. The thickness of the non-magnetic layer is chosen so that the magnetic layers are sufficiently far apart for exchange effects to be negligible but are close enough to be within the mean free path of conduction electrons in the material. If the two magnetic layers are magnetized in opposite directions and a current is passed through them along the direction of magnetization, half the electrons in each layer will be subject to increased scattering while half will be unaffected (to a first approximation). Furthermore, only the unaffected electrons will have mean free paths long enough for them to have a high probability of crossing the non magnetic layer. Once these electrons have crossed the non-magnetic layer, they are immediately subject to increased scattering, thereby becoming unlikely to return to their original side, the overall result being a significant increase in the resistance of the entire structure.\nIn order to make use of the GMR effect, the direction of magnetization of one of the layers must be permanently fixed, or pinned. Pinning is achieved by first magnetizing the layer (by depositing and/or annealing it in the presence of a magnetic field) and then permanently maintaining the magnetization by over coating with a layer of antiferromagnetic (AFM) material. The other layer, by contrast, is a xe2x80x9cfree layerxe2x80x9d whose direction of magnetization can be readily changed by an external field (such as that associated with a bit at the surface of a magnetic disk). Structures in which the pinned layer is at the top are referred to as top spin valves. Similarly, in a bottom spin valve structure the pinned layer is at the bottom.\nAlthough not directly connected to the GMR effect, an important feature of spin valve structures is a pair of longitudinal bias stripes that are permanently magnetized in a direction parallel to the long dimension of the device. Their purpose is to prevent the formation of multiple magnetic domains in the free layer portion of the GMR sensor, particularly near its ends. Thus longitudinal bias is responsible for the stability of a spin-valve recording head. It is usually achieved by an abutted-type junction formed during hard bias and lead deposition. Longitudinal bias layers may be implemented as permanent magnets by using a hard magnetic material or they may be compounded of a magnetically soft material whose magnetization is maintained by means of a contiguous antiferromagnetic layer.\nAs the dimensions of spin valve heads continue to be reduced, a number of difficulties in the associated fabrication processes are being encountered. When permanent magnets are used to supply contiguous junction stabilization, the junction region is always a major source of noise and instability due to the uncertain shape and coercivity reduction of the hard bias film at the junction area. It is known that the longitudinal bias permanent magnet is not as hard at the junction region as it is in the bulk material because of its long tapered tail. An external magnetic field can cause irreversible reversal of the hard magnet, thereby causing a hysteric response by the head. This coercivity reduction becomes more severe as the hard magnet becomes thinner.\nIn U.S. Pat. No. 5,664,316 it was shown that a ferromagnetic/antiferromagnetic coupled layer could be used to replace a permanent magnet. It was also claimed that when a ferromagnetic/antiferromagnetic coupled layer is used, magnetic instability in the junction area will be reduced. In U.S. Pat. No. 5,528,440 it was proposed to apply this structure to provide longitudinal bias in a spin valve head.\nIn both of the above-cited patents, the ferromagnetic/antiferromagnetic bi-layer was deposited directly after the removal of the side region of the sensor stack. This approach requires very precise control of the end point of the etching process. As an example, consider a bottom spin valve structure in which the GMR stack is ordered from top to bottom as: free/non-magnetic/AP1/Ru/AP2/AFM, where AP1 and AP2 are two ferromagnetic layers magnetized to be anti-parallel to one another, AP2 being the layer closest to the AFM.\nWhen such a structure is being made, etching (to form the GMR pedestal) cannot be allowed to proceed beyond the nonmagnetic layer. Etching down to AP1, AP2, or the AFM layer will result in head instability due to the transverse exchange coupling field from the AP or the AFM layer. As sensor thickness is reduced even further, this scheme becomes more and more difficult to control due to the criticality of the end point.\nAnother approach to avoiding the above-discussed problem has been to over etch so all GMR stack layers get removed along the pedestal\"\"s sides. However, this method introduces a new problem viz. an increase in the incidence of shorts between the sensor to the bottom shield since some of the lower dielectric layer must also get removed.\nA routine search of the prior art was performed with the following references of interest being found.\nIn U.S. Pat. No. 6,185,078, Lin et al. disclose using a layer of nickel oxide as a pinning layer for a NiFe bias layer. Since the NiO is an insulator, any insulation lost as a result of removing too much of the lower dielectric gets replaced. A glue layer of Ta is used under the lead layer. U.S. Pat. No. 5,664,316 (Chen et al.), U.S. Pat. No. 5,705,973 (Yuan et al.) are all related patents. U.S. Pat. No. 5,856,897 (Mauri) shows a stabilization layer under the lead layer.\nIt has been an object of the present invention to provide a bottom spin valve structure with longitudinal bias of improved stability.\nAnother object has been to provide a process for manufacturing said structure.\nA further object has been that said process allow substantial latitude for end point control during the etching of the GMR pedestal.\nThese objects have been achieved by inserting a decoupling layer between the antiferromagnetic layer that is used to stabilize the pinned layer of the spin valve itself and the soft ferromagnetic layer that is used for longitudinal biasing. This soft ferromagnetic layer is pinned by a second antiferromagnetic layer deposited on it on its far side away from the first antiferromagnetic layer. The presence of the decoupling layer ensures that the magnetization of the soft layer is determined only by the second antiferromagnetic layer."}
{"text": "Refrigerant vapor compression systems are well known in the art. Air conditioners and heat pumps employing refrigerant vapor compression cycles are commonly used for cooling or cooling/heating air supplied to a climate controlled comfort zone within a residence, office building, hospital, school, restaurant or other facility. Refrigerant vapor compression systems are also commonly used for cooling air to provide a refrigerated environment for food items and beverage products within display cases in supermarkets, convenience stores, groceries, cafeterias, restaurants and other food service establishments.\nConventionally, these refrigerant vapor compression systems include a compressor, a condenser, an expansion device, and an evaporator connected in refrigerant flow communication. The aforementioned basic refrigerant system components are interconnected by refrigerant lines in a closed refrigerant circuit and arranged in accord with the vapor compression cycle employed. An expansion device, commonly an expansion valve or a fixed-bore metering device, such as an orifice or a capillary tube, is disposed in the refrigerant line at a location in the refrigerant circuit upstream with respect to refrigerant flow of the evaporator and downstream of the condenser. The expansion device operates to expand the liquid refrigerant passing through the refrigerant line running from the condenser to the evaporator to a lower pressure and temperature. In doing so, a portion of the liquid refrigerant traversing the expansion device expands to vapor. As a result, in conventional refrigerant vapor compression systems of this type, the refrigerant flow entering the evaporator constitutes a two-phase mixture. The particular percentages of liquid refrigerant and vapor refrigerant depend upon the particular expansion device employed and the refrigerant in use, for example R-12, R-22, R-134a, R404A, R-410A, R-407C, ammonia, carbon dioxide or other compressible fluid.\nIn some refrigerant vapor compression systems, the evaporator is a parallel tube heat exchanger. Such heat exchangers have a plurality of parallel refrigerant flow paths therethrough provided by a plurality of tubes extending in parallel relationship between an inlet header and an outlet header. The inlet header receives the refrigerant flow from the refrigerant circuit and distributes the refrigerant flow amongst the plurality of flow paths through the heat exchanger. The outlet header serves to collect the refrigerant flow as it leaves the respective flow paths and to direct the collected flow back to the refrigerant line for return to the compressor in a single pass heat exchanger or through an additional bank of heat exchange tubes in a multi-pass heat exchanger.\nHistorically, parallel tube heat exchangers used in such refrigerant vapor compression systems have used round tubes, typically having a diameter of ⅜ inch or 7 millimeters. More recently, flat, rectangular dimension, multi-channel tubes are being used in heat exchangers for refrigerant vapor compression systems. Each multi-channel tube has a plurality of flow channels extending longitudinally in parallel relationship the length of the tube, each channel providing a small flow area refrigerant flow path. Thus, a heat exchanger with multi-channel tubes extending in parallel relationship between the inlet and outlet headers of the heat exchanger will have a relatively large number of small flow area refrigerant flow paths extending between the two headers. In contrast, a parallel tube heat exchanger with conventional round tubes will have a relatively small number of large flow area flow paths extending between the inlet and outlet headers.\nA problem associated with heat exchangers having flat, rectangular tubes extending between an inlet header and an outer header versus heat exchangers having round tubes is the connection of the inlet ends of the tubes to the inlet header. Conventionally, the inlet header is an axially elongated cylinder of circular cross-section provided with a plurality of rectangular slots cut in its wall at axially spaced intervals along the length of the header. Each slot is adapted to receive the inlet end of one of the flat, rectangular heat exchange tubes with the inlets to the various flow channels open to the chamber of the header, whereby fluid within the chamber of the inlet header may flow into the multiple flow channels of the various heat exchange tubes opening into the chamber. As the flat, rectangular heat exchange tubes have a lateral dimension significantly greater than the diameter of conventional round tubes, the diameters of the round cylindrical headers associated with conventional flat tube heat exchangers are significantly greater than the diameters of headers associated with round tube heat exchangers for a comparable volumetric fluid flow rate.\nNon-uniform distribution, also referred to as maldistibution, of two-phase refrigerant flow is common problem in parallel tube heat exchangers which adversely impacts heat exchanger efficiency. Two-phase maldistribution problems are caused by the difference in density of the vapor phase refrigerant and the liquid phase refrigerant present in the inlet header due to the expansion of the refrigerant as it traversed the upstream expansion device.\nOne solution to control refrigeration flow distribution through parallel tubes in an evaporative heat exchanger is disclosed in U.S. Pat. No. 6,502,413, Repice et al. In the refrigerant vapor compression system disclosed therein, the high pressure liquid refrigerant from the condenser is partially expanded in a conventional in-line expansion value upstream of the heat exchanger inlet header to a lower pressure, liquid refrigerant. A restriction, such as a simple narrowing in the tube or an internal orifice plate disposed within the tube, is provided in each tube connected to the inlet header downstream of the tube inlet to complete expansion to a low pressure, liquid/vapor refrigerant mixture after entering the tube.\nAnother solution to control refrigeration flow distribution through parallel tubes in an evaporative heat exchanger is disclosed in Japanese Patent No. JP4080575, Kanzaki et al. In the refrigerant vapor compression system disclosed therein, the high pressure liquid refrigerant from the condenser is also partially expanded in a conventional in-line expansion value to a lower pressure, liquid refrigerant upstream of a distribution chamber of the heat exchanger. A plate having a plurality of orifices therein extends across the chamber. The lower pressure liquid refrigerant expands as it passes through the orifices to a low pressure liquid/vapor mixture downstream of the plate and upstream of the inlets to the respective tubes opening to the chamber.\nJapanese Patent No. JP2002022313, Yasushi, discloses a parallel tube heat exchanger wherein refrigerant is supplied to the header through an inlet tube that extends along the axis of the header to terminate short of the end the header whereby the two phase refrigerant flow does not separate as it passes from the inlet tube into an annular channel between the outer surface of the inlet tube and the inside surface of the header. The two phase refrigerant flow thence passes into each of the tubes opening to the annular channel.\nObtaining uniform refrigerant flow distribution amongst the relatively large number of small flow area refrigerant flow paths is even more difficult than it is in conventional round tube heat exchangers and can significantly reduce heat exchanger efficiency. Two-phase maldistribution problems may be exacerbated in inlet headers associated with conventional flat tube heat exchangers due to the lower fluid flow velocities attendant to the larger diameter of such headers. At lower fluid flow velocities, the vapor phase fluid more readily separates from the liquid phase fluid. Thus, rather than being a relatively uniform mixture of vapor phase and liquid phase fluid, the flow within the inlet header will be stratified to a greater degree with a vapor phase component separated from the liquid phase component. As a consequence, the fluid mixture will undesirably be non-uniformly distributed amongst the various tubes, with each tube receiving differing mixtures of vapor phase and liquid phase fluid.\nIn U.S. Pat. No. 6,688,138, DiFlora discloses a parallel, flat tube heat exchanger having an inlet header formed of an elongated outer cylinder and an elongated inner cylinder disposed eccentrically within the outer cylinder thereby defining a fluid chamber between the inner and outer cylinders. The inlet end of each of the flat, rectangular heat exchange tubes extend through the wall of the outer cylinder to open into the fluid chamber defined between the inner and outer cylinders.\nJapanese Patent No. 6241682, Massaki et al., discloses a parallel flow tube heat exchanger for a heat pump wherein the inlet end of each flat, multi-channel tube connecting to the inlet header is crushed to form a partial throttle restriction in each tube just downstream of the tube inlet. Japanese Patent No. JP8233409, Hiroaki et al., discloses a parallel flow tube heat exchanger wherein a plurality of flat, multi-channel tubes connect between a pair of headers, each of which has an interior which decreases in flow area in the direction of refrigerant flow as a means to uniformly distribute refrigerant to the respective tubes."}
{"text": "1. Field of the Invention\nThe invention relates to a method for absolute calibration of an interferometer with the aid of a spherical output wave.\nThe invention relates to a method for determining the measurement accuracy of a recording device for a diffractive optical element.\n2. Description of the Related Art\nTwo methods are used in practice for the absolute determination of the measurement error of interferometers for the purpose of checking the fit of spherical surfaces, specifically the classic three-position test and the rotating-disc method. The main disadvantage of the known methods resides in the low number of absolute calibration methods for interferometers with a spherical output wave. Each method is subject to errors, and in order to determine these more closely it would be very advantageous to have additional absolute calibration methods for comparison."}
{"text": "1. Field of the Invention\nThe invention concerns a tool head for use in a machine tool, the machine tool having base body, a tool shank adapted to be coupled to a rotating machine spindle projecting axially beyond the base body and two cutting insert receptacles spaced apart in the circumferential and/or in the axial direction for receiving respectively one indexable cutting insert, of which the active main cutting surfaces exhibit deferring adjustment angles relative to the axis of the base body.\n2. Description of the Related Art\nIn known machine tools of this type (DE-A-196 05 156), which are designed for finishing valve seat rings and tappet guide bores in cylinder heads, three cutting inserts are provided arranged distributed about the circumference of the base body, wherein one is provided for producing the tight seat bevel for the valve and the other two are provided for producing an inlet and an and outlet protective bevel. The tight seat bevel is associated with a high precision requirement. For producing the three beveled areas, until now three different cutting inserts were employed, which are secured to the tool head by clamping in different positions and orientations. Each of these cutting inserts has a predetermined bevel angle for the associated purpose. In the case of wear then, during tool change-out, care must be taken to ensure that on the appropriate insert receptacle respectively the correct type of replacement insert is introduced. Disadvantages associated therewith include the difficulty of manipulation during change-out of cutting inserts and the high cost of cutting inserts."}
{"text": "Lighting applications typically use incandescent or gas-filled bulbs. Such bulbs typically do not have long operating lifetimes and thus require frequent replacement. Gas-filled tubes, such as fluorescent or neon tubes, may have longer lifetimes, but operate using high voltages and are relatively expensive. Further, both bulbs and gas-filled tubes consume substantial amounts of energy.\nA light emitting diode (LED) is a device that emits light upon the recombination of electrons and holes in an active layer of the LED. An LED typically includes a chip of semiconducting material doped with impurities to create a p-n junction. Current flows from the p-side, or anode, to the n-side, or cathode. Charge-carriers—electrons and holes—flow into the p-n junction from electrodes with different voltages. When an electron meets a hole, the electron recombines with the hole in a process that may result in the radiative emission of energy in the form of one or more photons (hν). The photons, or light, are transmitted out of the LED and employed for use in various applications, such as, for example, lighting applications and electronics applications.\nLED's, in contrast to incandescent or gas-filled bulbs, are relatively inexpensive, operate at low voltages, and have long operating lifetimes. Additionally, LED's consume relatively little power and are compact. These attributes make LED's particularly desirable and well suited for many applications.\nDespite the advantages of LED's, there are limitations associated with such devices. Such limitations include materials limitations, which may limit the efficiency of LED's; structural limitations, which may limit transmission of light generated by an LED out of the device; and manufacturing limitations, which may lead to high processing costs. Accordingly, there is a need for improved LED's and methods for manufacturing LED's."}
{"text": "Integration of a microphone in an integrated circuit may require formation of a back cavity to attain desired levels of microphone sensitivity. Formation of a back cavity with a sufficient volume may increase fabrication cost and complexity of the integrated circuit."}
{"text": "This invention relates to sensing of cardiac activity via electrocardiographic (ECG) data. In particular, this invention relates to a new and improved method of gathering ECG data and to the transmission of the collected data to a remote location for analysis. Alternatively, such ECG data may also be transmitted to a co-implanted device to provide data for the delivery of various therapies including pacing, cardioversion and defibrillation therapies, drug delivery, to effect capture detection and automatic stimulation threshold adaptation, to record PMT (pacemaker mediated tachycardia) episodes, to measure refractory periods, to set timing windows for anti-tachycardia pacing therapies, and to adjust the pacing rate to physiologic demand.\nThe electrocardiogram (ECG) is commonly used in medicine to determine the status of the electrical conduction system of the human heart. As currently practiced, the ECG recording device is commonly attached to the patient via ECG leads connected to pads arrayed on the patient\"\"s body so as to achieve a recording that displays the cardiac waveforms in any one of 12 possible vectors.\nThe detection and discrimination between various arrhythmic episodes in order to trigger the delivery of an appropriate therapy is of considerable interest. Prescription for implantation and programming of the implanted device are based on the analysis of the PQRST electrocardiogram (ECG) that currently requires externally attached electrodes and the electrogram (EGM) that requires implanted pacing leads. The EGM waveforms are usually separated for such analysis into the P-wave and R-wave in systems that are designed to detect the depolarization of the atrium and ventricle respectively. Such systems employ detection of the occurrence of the P-wave and R-wave, analysis of the rate, regularity, and onset of variations in the rate of recurrence of the P-wave and R-wave, the morphology of the P-wave and R-wave and the direction of propagation of the depolarization represented by the P-wave and R-wave in the heart. The detection, analysis and storage of such EGM data within implanted medical devices are well known in the art. Acquisition and use of ECG tracing(s), on the other hand, has generally been limited to the use of an external ECG recording machine attached to the patient via surface electrodes of one sort or another.\nRecently, however, an alternative method of collecting ECG tracings from a set of subcutaneous electrodes, or a subcutaneous electrode array (SEA), has been described in the following patents or patent applications. U.S. Pat. No. 5,331,966 to Bennett, incorporated herein by reference, discloses a method and apparatus for providing an enhanced capability of detecting and gathering electrical cardiac signals via an array of relatively closely spaced subcutaneous electrodes (located on the body of an implanted device).\nMore recently, a patent application entitled xe2x80x9cSurround Shroud Connector and Electrode Housings for a Subcutaneous Electrode Array and Leadless ECGsxe2x80x9d, by Ceballos, et al., filed on Oct. 26, 2000, Ser. No. 09/697,438, incorporated herein by reference in its totality, discloses an alternate method and apparatus for detecting electrical cardiac signals via an array of subcutaneous electrodes located on a shroud circumferentially placed on the perimeter of an implanted pacemaker. Similarly, a patent application entitled xe2x80x9cSubcutaneous Electrode for Sensing Electrical Signals of the Heartxe2x80x9d, by Brabec et al, filed Oct. 31, 2000, Ser. No. 09/703,152, incorporated herein by reference in its totality, discloses the use of a spiral electrode using in conjunction with the shroud described in the Ceballos et al disclosure. In addition, two applications, entitled xe2x80x9cMultilayer Ceramic Electrodes for Sensing Cardiac Depolarization Signalsxe2x80x9d, by Guck et al, filed Oct. 25, 2000, Ser. No. 09/696,365 and P-8787 xe2x80x9cThin Film electrodes for Sensing Cardiac Depolarization Signalsxe2x80x9d by Guck and Donders, filed Dec. 13, 2000, Ser. No. as yet unknown, both incorporated herein by reference in their totality, disclosed the use of sensing electrodes placed into recesses incorporated along and into the peripheral edge of the implantable pacemaker. Finally, the submission entitled, xe2x80x9cSubcutaneous Electrode Array Virtual ECG Leadxe2x80x9d by Panken and Reinke, filed Nov. 22, 2000, Ser. No. 09/721,275, also incorporated by reference to its entirety, describes the algorithm used by the implanted device that compiles the ECG from any two subcutaneous electrodes found in the SEA.\nAs the functional sophistication and complexity of implantable medical device systems increased over the years, it became increasingly more important for such systems to include a system for facilitating communication between one implanted device and another implanted device and/or an external device, for example, a programming console, monitoring system, or the like. For diagnostic purposes, it is desirable that the implanted device be able to communicate information regarding the device\"\"s operational status and the patient\"\"s condition to the physician or clinician. State of the art implantable devices are available which can even transmit a digitized electrical signal to display electrical cardiac activity (e.g., an ECG, EGM, or the like), programmed pacing parameters, therapy status, etc. for storage and/or analysis by an external device.\nTo diagnose and measure cardiac events, the cardiologist has several tools from which to choose. Such tools include twelve-lead electrocardiograms, exercise stress electrocardiograms, Holter monitoring, radioisotope imaging, coronary angiography, myocardial biopsy, and blood serum enzyme tests. In addition, a record of the programmed parameters, the status of atrial and ventricular auto-capture, the atrial and ventricular sensing threshold, and other such pacing system parameters are required for an adequate diagnosis. Such ECG tracings and printed reports are placed into the patient\"\"s records and used for comparison to more recent tracings and reports. It must be noted, however, that whenever an ECG recording is required (whether through a direct connection to an ECG recording device or to a pacemaker programmer), external electrodes and leads must be used. Moreover, the acquisition of programmed data and pacing therapy status requires the use of a programming system that is in immediate proximity to the patient.\nUnfortunately, surface ECG electrodes have some serious drawbacks. For example, electrocardiogram analysis performed using existing external or body surface ECG systems is likely to be limited by mechanical problems and poor signal quality. Electrodes attached externally to the body are a major source of signal quality problems and analysis errors because of susceptibility to interference such as muscle noise, power line interference, high frequency communication equipment interference, and baseline shift from respiration or motion. Signal degradation also occurs due to contact problems, ECG waveform artifacts, and patient discomfort. Externally attached electrodes are subject to motion artifacts from positional changes and the relative displacement between the skin and the electrodes. Furthermore, external electrodes require special skin preparation to ensure adequate electrical contact. Such preparation, along with positioning the electrode and attachment of the ECG lead to the electrode needlessly prolongs the pacemaker follow-up session. One possible approach is to equip the implanted pacemaker with the ability to detect cardiac signals, transform them into a tracing, and transmit them to a remote receiving device such as a programmer, transtelephonic pacemaker monitoring system, or a central database system. Once captured, these ECG tracing could be available for use by medical personnel for analysis and diagnosis.\nSimilarly, it is often difficult if not impossible to have a programming unit nearby to gain a visual or printed report or all the various programmed parameters and the status of the associated pacing therapies. For example, when a physician visits a hospitalized patient, it is now mandatory that he or she carry a programmer to the patient\"\"s bedside to acquire a full status of the pacemaker\"\"s operation. If, for one reason or another, the physician has no programmer immediately available, he or she must take the time to retrace steps and find one to bring back to the patient\"\"s bedside. It would certainly be advantageous if that patient\"\"s ECG and pacemaker status could be acquired remotely prior to the physician\"\"s approach to the bedside. In addition, when a pacemaker patient visits his or her physician, there is often a period of time that the patient sits in the waiting room. This time could be advantageously to capture the patient\"\"s ECG data as well as the pacemaker\"\"s programmed status via a transmission to a programmer located in the pacemaker technician\"\"s office. Then, these reports could be printed and made available to the physician for study prior to bringing the patient into the pacemaker clinic. The patient, especially the female patient, would feel more comfortable, not having to disrobe to complete the pacemaker follow-up session.\nAs pacing and other medical technology have advanced, it is often necessary to co-implant other implantable devices such as a cardioverter-defibrillator, drug delivery system, a pressure sensor, and so on. Each of these devices may require ECG data upon which to make therapy and diagnostic decisions. Thus the ECG data already available in the implanted pacemaker may be transmitted intracorporeally to provide control signals for the delivery of various therapies including pacing, cardioversion and defibrillation therapies as well as the delivery of antiarrhythmic drugs, and, in the pacing context, to effect capture detection and automatic stimulation threshold adaptation, recording of PMT episodes, measurement of refractory periods in order to set timing windows for antitachycardia pacing therapies, and as a control signal for use in adjusting pacing rate to physiologic demand. All these devices would benefit from enhanced capabilities of discriminating cardiac arrhythmias and storing the associated data.\nAs used herein, Capture is defined as an evoked cardiac response to a pacemaker output or stimulation pulse. In a dual chamber pacemaker, for example, a stimulation pulse can be applied to either the atrium or the ventricle during the appropriate portion of a cardiac cycle. The minimum output pulse energy required to capture and thus evoke a cardiac depolarization is referred to as the stimulation threshold. This threshold generally varies in accordance with the well known strength-duration curves, wherein the amplitude of a stimulation threshold current pulse and its duration are inversely proportional.\nA number of factors can influence changes in the stimulation threshold following implantation of the pacemaker and pacing lead. These include: (1) changes in position of the pacing electrode relative to the cardiac tissue; (2) long-term biologic changes in cardiac tissue closely adjacent the electrode, such as fibrotic tissue growth; (3) changes in the patient\"\"s sensitivity to stimulation as a function of periodically fluctuating conditions, even on a daily basis, due to various causes such as diet, exercise, administered drugs, electrolyte changes, etc.; and (4) gradual changes in pacemaker/lead performance due to various causes such as battery depletion, component aging, etc.\nCapture Threshold Detection and Adjustable Output Pulse Energy require that battery power be conserved to extend the pacemaker\"\"s useful life. It is therefore desirable to achieve capture at the lowest possible settings (amplitude and pulse width) for the output pulse. With the advancement of programmable pacemakers, it became common to initially program an output pulse energy setting which includes a safety margin somewhat above that required to produce capture. These programmable pacemakers include programmable options that permit the physician to select output pulse settings known to be sufficient to capture the heart but which are below the maximum obtainable output voltage of the pacemaker. The attending physician usually determines such output pulse adjustments during an office visit with the use of an external programmer and an electrocardiogram (ECG) monitor. At this time, the physician assesses capture by means of an ECG, while at the same time the pacemaker is providing a sequence of temporarily programmed-in stimulation pulses with decreasing pulse voltages in a system of the type described in\nU.S. Pat. No. 4,250,884 to Hartlaub, et al. For example, capture detection of the ventricle is confirmed by the presence of the evoked QRS complex or R-wave, and capture detection of the atrium is confirmed by the presence of the evoked P-wave. Loss of capture can be directly observed and correlated to the pulse energy at which capture is lost.\nSince the late 1960\"\"s, self-adaptive pacemakers have been developed that have the capability of automatically adjusting the voltage of the pacing pulse as appropriate to accommodate changes in stimulation threshold.\nU.S. Pat. No. 3,757,792 to Mulier et al, for example, relates to an early pacemaker that provides for a decreased battery drain by sensing each driven heart beat (i.e., R-wave) and providing for a decrease in energy for each succeeding output pulse until such time as loss of capture is detected. Following a detected loss of capture, the next succeeding output pulse is increased in voltage by an amount to be safely over the threshold hysteresis level. U.S. Pat. No. 3,949,758 to Jirak (incorporated herein by reference) relates to a similar threshold-seeking pacemaker with automatically adjusted voltage levels for output pulses in response to detected loss of capture (i.e., absence of R-wave), and describes separate sensing and pacing electrodes, which are each utilized in unipolar fashion with a third common electrode having a comparatively larger dimension, to reduce residual polarization problems.\nU.S. Pat. No. 3,977,411 to Hughes, Jr. et al shows a pacemaker having separate sensing and pacing electrodes that are each utilized in unipolar fashion. The sensing electrode comprises a ring electrode having a relatively large surface area (i.e., between 75 to 200 mm.sup.2) for improved sensing of cardiac activity (R-waves), and is spaced along the pacing lead approximately 5 to 50 mm from the distally-located tip electrode used for pacing.\nU.S. Pat. No. 3,920,024 to Bowers shows a pacemaker with a threshold tracking capability that dynamically measures the stimulation threshold by monitoring the presence or absence of an evoked response (R-wave). If no R-wave is sensed within a post-stimulus interval (e.g., 20 to 30 ms post-stimulus), the pacemaker delivers a closely spaced backup pulse (e.g., 40 to 50 ms post-stimulus) at increased amplitude and pulse width to ensure an evoked response.\nU.S. Pat. Nos. 4,759,366 and 4,858,610 to Callaghan, et al, incorporated herein by reference, relate to evoked response detector circuits which also employ fast recharge in at least one separate sensing electrode in either unipolar or bipolar electrode configurations in either or both the atrium and ventricle. The cardiac pacing systems function as unipolar and bipolar systems at different steps in the operating cycle. In the \"\"610 patent, a separate electrode on the connector block of the pacemaker can is suggested for use as the reference electrode anode rather than the metal case itself if the case is employed as the reference electrode for the delivery of the stimulation pulse. In the \"\"366 patent, the detected evoked response is used in an algorithm for adjusting the pacing rate.\nU.S. Pat. No. 4,310,000 to Lindemans and U.S. Pat. Nos. 4,729,376 and 4,674,508 to DeCote, incorporated herein by reference, also disclose the use of a separate passive sensing reference electrode mounted on the pacemaker connector block or otherwise insulated from the pacemaker case in order to provide a sensing reference electrode which is not part of the stimulation reference electrode and thus does not have residual after-potentials at its surface following delivery of a stimulation pulse. The DeCote \"\"376 and \"\"508 patents also set forth stimulation threshold testing algorithms for adjusting the pacing pulse energy.\nThus, considerable effort has been expended in providing electrode systems, fast recharge circuitry and separate sense amplifier systems for avoiding after-potentials and providing capture detection and stimulation threshold tracking.\nDistinguishing malignant tachyarrhythmias from sinus tachycardias and detecting pacemaker mediated tachycardias is similarly limited by the available electrode configurations employed in single and dual chamber pacing systems, implantable drug dispensers and pacemaker-cardioverter-defibrillator systems as described above. In the context of discriminating spontaneously occurring tachyarrhythmias from sinus tachycardia, attempts have been made to employ both atrial and ventricular electrode systems in order to determine whether the tachycardia is sinus in origin or reflects a retrogradely conducted abnormal ventricular rhythm. For example, previous art has placed multiple electrodes on atrial and ventricular leads and to sense the direction of travel of a depolarization wave front as shown for example in U.S. Pat. No. 4,712,554 to Garson, Jr.\nIn addition, other art has shown that pacemakers that operated in the DDD or related modes can, under certain circumstances, sustain a dangerous tachycardia with ventricular pacing at the upper rate limit, a pacemaker mediated tachycardia (PMT). A PMT may be initiated when AV synchrony is lost, leaving the AV node conductive tissue able to conduct retrograde electrical signals from the ventricle to the atrium which in turn cause an atrial depolarization. The sensing of the resulting atrial depolarization by the atrial sense amplifier in turn causes the ventricular pulse generator to emit a ventricular pacing pulse after the AV time period times out. The cycle may repeat itself if the ventricular pace event is conducted to the atrium where it again causes an atrial depolarization which is picked up by the atrial sense amplifier. This repetitive high rate stimulation may be sustained indefinitely by the pacemaker causing discomfort to the patient or possibly inducing more threatening arrhythmias.\nVarious techniques have been implemented to minimize the impact of PMTs, but these techniques usually sacrifice flexibility of the DDD system. U.S. Pat. No. 4,967,746 to Vandegriff sets forth a number of techniques which have been employed to alleviate PMTs.\nECG/EGM Vector Analysisxe2x80x94xe2x80x94U.S. Pat. No. 4,310,000 issued to Lindemans suggests various modifications to the passive sensing reference electrode depicted in its drawings, including the incorporation of more than one passive sensing reference electrode provided on or adjacent to the metallic can, positioned as deemed necessary for best sensing, and connected to one or more sense amplifiers. No specific use of the additional passive sensing reference electrodes is suggested, although the single passive sensing reference electrode is suggested for use with a sense amplifier to detect both capture and spontaneous atrial or ventricular electrical events in a dual chamber pacing system.\nIt is known in the art to monitor electrical activity of the human heart for diagnostic and related medical purposes. U.S. Pat. No. 4,023,565 issued to Ohlsson describes circuitry for recording EKG signals from multiple lead inputs. Similarly, U.S. Pat. No. 4,263,919 issued to Levin, U.S. Pat. No. 4,170,227 issued to Feldman, et al, and\nU.S. Pat. No. 4,593,702 issued to Kepski, et al, describe multiple electrode systems which combine surface EKG signals for artifact rejection.\nThe primary use for multiple electrode systems in the prior art appears to be vector cardiography from EKG signals taken from multiple chest and limb electrodes. This is a technique whereby the direction of depolarization of the heart is monitored, as well as the amplitude. U.S. Pat. No. 4,121,576 issued to Greensite discusses such a system.\nIn addition, U.S. Pat. No. 4,136,690 issued to Anderson, et al, shows a vector cardiographic system used for arrhythmia analysis. Similar techniques are described in xe2x80x9cRhythm Analysis Using Vector Cardiograms,xe2x80x9d Transactions on Biomedical Engineering, Vol. BME-32, No. 2, February 1985, by Reddy, et al, European Pat. No. 0 086 429 issued to Sanz and U.S. Pat. No. 4,216,780 issued to Rubel, et al.\nVarious systems have additionally been proposed for measuring the orthogonal ventricular or atrial electrogram from multi-electrode lead systems placed endocardially within the patient\"\"s atrium and/or ventricle. Such orthogonal endocardial EGM systems are depicted in U.S. Pat. No. 4,365,639, issued to Goldreyer, and U.S. Pat. Nos. 4,630,611 and 4,754,753 issued to King. In addition, orthogonal ventricular electrogram sensing employing endocardial, multi-electrode lead systems and associated circuitry are disclosed in two articles by Goldreyer, et al, entitled xe2x80x9cOrthogonal Electrogram Sensing,xe2x80x9d PACE, Vol. 6, pp. 464-469, March-April 1983, Part II, and xe2x80x9cOrthogonal Ventricular Electrogram Sensing,xe2x80x9d PACE, Vol. 6, pp. 761-768, July-August 1983. In the Goldreyer patent and in these papers, it is suggested that the orthogonal electrodes be employed to detect paced events and provide capture verification as well as to facilitate the discrimination of P-waves from QRS complexes. Other articles by Goldreyer, et al., appear in the literature, including those listed in the bibliographies to these two papers.\nThe aforementioned King U.S. Pat. Nos. 4,630,611 and 4,754,753 describe X, Y and Z orthogonally displaced electrodes on the body of the endocardial pacing lead and circuitry for developing a composite EGM vector signal in order to detect changes in the vector over time and discriminate normal sinus rhythm from tachyarrhythmias.\nFinally, U.S. patent application Ser. No. 611,901 entitled xe2x80x9cMulti-Vector Pacing Artifact Detector,xe2x80x9d filed Nov. 9, 1990, and assigned to the assignee of the present application, sets forth a system for detecting the artificial pacing artifact in patients having artificially paced myocardial contractions in an external monitor employing three standard EKG leads with chest or limb electrodes.\nThere are a number of tasks that are normally performed either prior to or during a pacemaker follow-up session. For example, the patient\"\"s chart must be pulled and made available for reference by the clinician or physician prior to seeing the patient. Then the patient must be prepped for collecting ECG tracings. Generally, such preparation requires disrobing from the waist up and attaching ECG electrodes at the required point on the chest and other points on the body. This is followed by placing the programming head over the site of the implanted pacemaker, collecting the ECG tracings, interrogating the pacemaker, printing the report, examining all the printouts, and then determining if the pacemaker\"\"s operation is optimal. If not, then programming change(s) are usually required. Another interrogation will be made and a final printed report (which may include new ECG tracings) is captured for inclusion in the patient\"\"s records. Other tasks may also be required and performed depending on the protocol to be followed by the institution or clinic to which the patient is assigned.\nOne might question why all these tasks must be performed. An examination of why ECG tracings are demanded will prove instructive. Magnet and non-magnet ECG tracings are normally required for clinic use and absolutely required for remote pacemaker follow-up. In fact, an institution may not be reimbursed for follow-up services unless such ECG tracings are obtained, as occurs in the U.S. The matter of reimbursement varies from country to country; but most institutions and clinics require such tracings for their records as a matter of medical completeness, if for no other reason. Thus, such ECG tracings have become a critical component of the pacemaker followup.\nAs mentioned, the collection of ECG tracings require the attachment of electrodes, such as chest electrodes for in-clinic follow-ups and wrist or fingertip electrodes for a remote, transtelephonic followup. The elderly patient at home often has difficulty attaching his or her wrist or fingertip electrode so that a suitable ECG tracing is recorded and transmitted. If the difficulty persists, the patient may then be asked to come into the clinic for his/her pacer check, a far more costly visit. The in-clinic collection of ECG tracings, while it does involve a trained technician, still requires a good deal of time for attaching the ECG electrodes, ensuring they are connected appropriately to the programmer, recording the ECG tracings, removal of the ECG electrodes and, finally, storing them and/or getting them out of the way so the remainder of the follow-up session may proceed.\nUnfortunately, the methods just described for collecting such ECG tracings within a clinic setting or during a remote followup have several issues. First, clinics are very sensitive to the duration of an in-clinic or remote follow-up session. The most time-consuming parts to the in-clinic procedure just described are the attachments and placement of the ECG electrodes and verification that the positioning of the electrodes results in satisfactory ECG tracings. By far, the most common issue with a remote follow-up session is the acquisition of a quality ECG tracing. It often requires several attempts by the patient to correctly position the ECG electrodes.\nSecondly, the quality of the ECG is highly dependent on a number of factors independent of electrode placement such as patient movement. While the potential for patient movement is greater during and remote follow-up session than during the in-clinic session, it does occur there as well.\nThirdly, it is often difficult to determine atrial activity on ECG tracings. This is particularly true of a remote follow-up. During an in-clinic session, however, if a digital ECG recorder is being used rather than the programmer\"\"s ECG capability, the clinician may be unable to determine whether the atria are being paced due to the pacing pulse being lost between digital processing times.\nFourthly, a trained technician, clinician, or physician must validate the quality of the ECG tracings during an in-clinic or remote follow-up session. Thus, medical personnel who are untrained in ECG interpretation cannot currently be involved in either an in-clinic or remote pacemaker follow-up session.\nFinally, current technology requires the use of a programmer to download all other data stored in the pacemaker\"\"s memory. It is currently impossible to download important and critical pacing data during a remote follow-up session. Similarly, an in-clinic follow-up session is needlessly prolonged by the current inability of implanted devices to telemeter data to a programmer or other devices that can receive transmitted pacemaker data and print it for study by the technician, clinician, or physician prior to actually seeing the patient.\nThese issues may be resolved by the use of telemetry of data to and from the implanted device. The FCC has adopted the following definition for describing wireless medical telemetry: xe2x80x9cthe measurement and recording of physiological parameters and other patient-related information via radiated bi- or unidirectional electromagnetic signals.xe2x80x9d Thus, in its broadest sense, telemetry can be defined as the art and science of conveying information from one location to another. With radio telemetry, radio signals are utilized to convey that information.\nTwo terms have been used to describe the type of data transmission envisaged by the current invention, Telehome and Telemedicine. Telehome care can be defined as providing monitoring (telemetry) and home health care services at a distance, using advanced telecommunications and information technology. Telemedicine is the use of telecommunications and information technology to provide clinical care at a distance. The definition of xe2x80x9cdistancexe2x80x9d may range from several yards, such as might occur within a clinic environment or hundreds of miles as occurs in rural areas. The availability of telemedicine services in far-flung rural areas can sometimes mean the difference between life and death for patients who must travel hundreds of miles to see a nurse or doctor. Wireless technologies can be particularly beneficial in these areas because developing wireless networks may be faster and cheaper than building a wire line infrastructure.\nSome emergency medical transport companies have already begun to outfit their ambulances with wireless telemedicine equipment. This equipment enables a paramedic to communicate with the emergency physician for an early assessment, well before the patient\"\"s arrival at the hospital. The telemedicine equipment can be as simple as a laptop computer with desktop videoconferencing capabilities that provide simultaneous two-way video, two-way voice, vital signs, cardiac and other data to a trauma center. This latter type of equipment is envisaged as a first step toward implementing the current invention, since wireless products such as handheld computers with Internet capabilities are becoming commonplace.\nFurther, medical telemetry equipment is used in hospitals and health care facilities to transmit patient measurement data to a nearby receiver, permitting greater patient mobility and increased comfort. Examples of medical telemetry equipment include heart, blood pressure and respiration monitors. The use of these devices allows patients to move around early in their recovery while still being monitored for adverse symptoms. With such devices, one health care worker can monitor several patients remotely, thus decreasing health care costs. In a similar manner, ECG and other data from an implanted device may be telemetered both intra- and extra-corporeally.\nMedical telemetry equipment is increasingly relied upon in hospitals to improve health care and reduce costs. Pacing patients who formerly required the monitoring and treatment capabilities that are available only in pacemaker clinics can be monitored and treated remotely in their homes. The number of pacemaker patients with chronic medical conditions is rising due to the growth in the elderly population. For these reasons, the need for monitoring these patients outside of the clinic setting is rapidly increasing, and this need can be fulfilled with medical telemetry equipment.\nThe Leadless Fully Automatic Pacemaker Followup (LFAPF) system of the present invention addresses all the issues described above and makes use of the type of telemetry described above. Briefly, this novel process involves the transmission of a signal to the pacemaker that may trigger one or any number of transmissions. ECG tracings can be transmitted for processing and printing, as well as the contents of the pacemaker\"\"s memory including pacemaker diagnostics such as the number and duration of mode switches. Such triggering would also include activation of a number of diagnostic operations, such as the measurement of the stimulation threshold and/or Capture Management."}
{"text": "There are many commercial and industrial environments where engines or other vibration inducing mechanisms are supported at relatively fixed structure with some type of vibration damping system interposed between the support structure and the vibrating mechanism. For example, internal combustion engine driven vehicles are provided with shock absorbing engine mounts which are intended to limit the transmission of engine vibrations into the vehicle body frame and vice versa.\nCertain prior art vehicle engine mounts comprise elastically deformable hard rubber cushions or the like for cushioning the transfer of vibrations between the engine and the vehicle frame. Such solid elastically deformable engine mounts do help isolate the vehicle body frame carrying the engine from engine induced mechanical vibrations and the engine from vehicle induced vibrations. However, such engine mounts suffer from disadvantages in that they are unable to attenuate the wide range of mechanical vibrations experienced in use on an automobile For example, an automobile is designed to operate under many speed, torque, acceleration and deceleration conditions, all of which create different mechanical vibration force, frequency and amplitude patterns. The above-mentioned conventional shock absorber type engine mounts are preset and passive in that they only passively react to vibration forces based on their preset elastic design characteristics. Such vibration isolation mounts can be designed to operate quite well over certain narrow mechanical vibration patterns. However, it is impractical, if not impossible, to design effective preset passive isolation vibration mounts which can attenuate vibrations over all operating ranges of the engine and vehicle.\nCertain other prior art vehicle engine mounts utilize a closed fluid shock absorber system for cushioning the transfer of vibrations between the engine and the vehicle frame. Since these fluid shock absorber mounts are closed systems, they act much like the elastically deformable spring or hard rubber cushion mounts discussed above. Further, in cases of large relative movement between the parts being supported, the shock attenuation substantially diminishes as the fluid pressure rises. Although certain shock absorber mounts provide for multiple fluid chambers and different flow paths depending upon the relative displacement of the parts being supported, such systems are very complex and costly to construct and are also ultimately limited by the total volume of fluid and fluid accommodating space in the overall closed system at an individual shock absorber mount.\nThe above-mentioned prior art arrangements thus permit an unacceptable level of transfer of vibration forces between the engine and the vehicle frame, and/or involve very complicated, expensive and space wasting constructions.\nThe mechanical vibration problems of internal combustion engine driven automotive vehicles discussed above are but one example of complex mechanical vibration problems that occur. Another example is a drive assembly support for the rotor blades of a helicopter Such helicopter rotor blade drive assemblies further compound the mechanical vibrations that should be attenuated. Prior attempts to attenuate vibrating force transfers between helicopter blade drive assemblies and the helicopter body or helicopter passenger compartment meet with similar difficulties as described above for the motor vehicles driving on the ground, due again to the wide spectrum of vibration force, frequency and amplitude which are experienced during the helicopter operation. Mechanical vibration isolation problems also occur in stationary engine environments such as driving engines for heating and air conditioning units mounted on commercial buildings or stationary engine electrical generators used at construction sites and at remote areas where electric power must be generated. In these installations, it is desired to minimize mechanical vibration force transfers between the engine and pumps or generators driven by same so as to minimize noise and minimize vibration induced mechanical failures\nAn exemplary environment where reduction of mechanical vibration transfer is especially desired involves so-called top floor \"penthouse\" spaces of large commercial buildings with rooftop heating, venting and air conditioning systems (HVAC systems). Due to vibration induced sound and actual vibration movement of the floors and walls caused by the rooftop HVAC systems, top floors of such buildings are not desirable prime office or living spaces. Again, prior attempts to passively mount the large generators, pumps, etc, for such HVAC systems can only imperfectly and inefficiently respond to the induced mechanical vibration and especially to changes in the vibration spectrum occurring during operations."}
{"text": "1. Field of the Invention\nThe present invention relates to a shift control system for an electric vehicle, for controlling the shifting of a transmission interposed between an electric motor and a driven wheel based on a shift map compiled with at least the accelerator opening degree and the vehicle speed used as parameters.\n2. Description of the Prior Art\nIn an electric vehicle using an electric motor connected to a battery as a traveling drive source, a technique is conventionally known in which a regenerative braking force is exhibited by the electric motor to convert kinetic energy of the moving vehicle into electric energy, and the electric energy is recovered in the battery (for example, see Japanese Patent Application Laid-open No. 191904/93).\nIn the technique described in Japanese Patent Application Laid-open No. 191904/93, in order to exhibit a regenerative braking force corresponding to an engine brake in an engine-mounted vehicle when the accelerator opening degree is smaller than a reference accelerator opening degree, the regenerative braking force at the reference accelerator opening degree is set at zero, and when the accelerator opening degree is at a value such that an accelerator pedal is in a fully closed state, a maximum regenerative braking force is set. The regenerative braking force is varied linearly between zero and the maximum regenerative braking force value.\nIn the known technique, however, if the vehicle speed is increased, the maximum regenerative braking force is decreased in a hyperbolic curve. For this reason, if the vehicle speed is increased, for example, when the vehicle travels down a downward slope, a speed-reducing force is abruptly weakened to produce an inconsistency of braking feeling."}
{"text": "The present disclosure relates to an image reading device, an image forming apparatus including the same, and an image reading method.\nGenerally, there is known a light source switching type image reading device that obtains color image data by sequentially lighting light sources (LED, or the like) of three colors corresponding to R (red), G (green), and B (blue) and by reading, at that time, light data received by a line sensor as R data, G data, and B data, respectively.\nIn the light source switching type image reading device, since image data of respective colors are read while a document and the line sensor are relatively moved in a sub-scanning direction at a constant speed, the reading positions of the respective colors in the document do not coincide and are shifted in the sub-scanning direction. As a result, for example, each color is sometimes different from that in the document at, for example, a boundary portion of the color. Under such circumstances, there is known an image reading device that is able to correct such a color shift by correcting the image data of the respective colors that are read."}
{"text": "Roller conveyors are commonly used in many industries to move packages or materials. In general, a conveyor roller comprises a plurality of individual rollers mounted between two parallel frames. The roller typically comprises a tube having a bearing in each end thereof which is rotatably journaled on a roller shaft that extends longitudinally through the tube. The roller shaft protrudes from each end of the roller tube and engages mounting holes in the frame.\nFrequently, the roller shaft not only serves to mount the roller, but also functions as a structural element of the conveyor frame. In such cases, the ends of the shaft are drilled and tapped. Bolts extending from the outside of the frame are threaded into the tapped holes in the roller shaft to secure the frame members and rollers together. This approach to the construction of conveyors makes assembly of the conveyor time consuming and difficult. Moreover, this method of constructing conveyors makes it difficult to replace a single roller.\nAnother method to mount rollers in a conveyor frame is to use a roller having a through-shaft which is spring-loaded. In this case, the frame members are held together independently by a series of cross members rather than through the rollers. The rollers generally are installed after the frame is assembled. To install the roller, the shaft is moved axially against the force of the spring and \"snaps\" into place when the shaft aligns with the mounting hole in the frame member. This method makes assembly of the conveyor much simpler and facilitates the replacement of rollers when they become worn or defective.\nSeveral attempts have been made in the past to eliminate through-shafts in the roller. In place of a single thru-shaft, separate stub shafts have been used to support opposing ends of the conveyor roller. Typically, these stub shafts are bolted to the conveyor frame before the frame is assembled. The conveyor roller is then inserted over the stub shaft as the frame is locked together to create a final assembly. This method of mounting rollers in conveyors also makes assembly difficult. Moreover, this method essentially precludes prefabrication of the conveyor frame.\nAccordingly, there is a need for a new method for mounting conveyor rollers in a frame which does not rely on a through-shaft extending through the roller, allows prefabrication of the conveyor frame, and facilitates ease of assembly."}
{"text": "Computers and computing systems have affected nearly every aspect of modern living. Computers are generally involved in work, recreation, healthcare, transportation, entertainment, household management, etc.\nComputers generally include hardware which allows the computer to output visual representations of data to a user. For example, computers may include a video card connected to a CRT monitor or LCD video screen which allows visual representations of objects including characters or other graphics to be displayed to the user. Similarly, many computers include the ability to interface with printers to provide visual representations of objects on a hard copy for a user.\nDepending on the nature of the visual representations and the medium used to convey the visual representations to the users, the visual representations output to a user may vary from an ideal output. For example, LCD video screens are typically only able to output data using a pixel resolution of about 110-120 pixels per inch. Outputting a visual representation of a character may result in the character including aliasing due to the size and resolution limitations of a screen.\nTo correct less than ideal representations, filters may be applied to an object. For character objects, filters are applied to the entire object. Thus, in the example above, an antialiasing filter such as gray scaling may be applied to an entire character to help soften the effects of the aliasing.\nThe subject matter claimed herein is not limited to embodiments that solve any disadvantages or that operate only in environments such as those described above. Rather, this background is only provided to illustrate one exemplary technology area where some embodiments described herein may be practiced."}
{"text": "Electronic devices such as a smartphone, a personal computer, and a tablet computer provide many useful functions to users through various applications. These electronic devices are being developed to additionally provide various types of information by various functions as well as a voice call function.\nAside from a simple voice call function and an Internet browsing function, users of electronic devices have recently demanded various entertainment functions.\nAs one of the entertainment functions, a function of displaying an element(s) (for example, a piano keyboard) of a user-intended musical instrument (for example, a piano) and allowing a user to play the musical instrument using the displayed element(s) is provided.\nHowever, since the elements of the musical instrument are displayed in a very limited space such as a display of an electronic device, it might be difficult for the user has difficulty in playing the musical instrument using the displayed element(s). Accordingly, the need exists for new techniques for simulating musical instruments."}
{"text": "This invention relates to a winding device for a vehicle safety belt having a housing, a sleeve rotatable in the housing and spring-biased to wind up the safety belt, a ratchet wheel with ratchet teeth arranged at one end of the sleeve, counter teeth disposed on the housing, a catch wheel with catch teeth arranged between the ratchet wheel and housing and being displaceable into a locking position of simultaneous engagement with the ratchet teeth and counter teeth, and tension limiting means connecting the ratchet wheel for rotation with the sleeve supporting the safety belt and consisting of a torsion bar extending coaxially through the sleeve and connected at one end thereof with the sleeve and at the other end with the ratchet wheel.\nDuring rapid deceleration of the vehicle, such as on impact with another object, the body of the passenger seated in the vehicle is thrown against the restraining safety belt, thereby subjecting the body to dangerously high forces. In safety belt winding devices, it is known that these forces may be dissipated through use of a torsion bar placed between the sleeve carrying the safety belt and the ratchet wheel which on impact becomes locked in engagement with the safety belt winding device housing, as shown in German Offenlegungsschrift No. 2,064,710, for example. When pre-determined, acceptable belt forces are exceeded, the torsion bar placed between the sleeve and ratchet wheel is twisted thereby absorbing energy.\nHowever, it has been found that when the torsion bar is twisted, it undergoes an increase in axial length. The lengthening of the torsion bar may result in disengagement of the ratchet wheel from its locking position due to the lengthening torsion bar's lifting the ratchet wheel away from the locking catch wheel interposed on the sleeve between the ratchet wheel and the safety belt housing. Such disengagement seriously impairs the operability and effectiveness of the safety belt winding device as a whole."}
{"text": "Pathologies of the gastro-intestinal (“GI”) system, the biliary tree, the vascular system and other body lumens and hollow organs are commonly treated through endoscopic procedures, many of which require active and/or prophylactic hemostasis to reduce internal bleeding. Tools for deploying hemostatic clips via endoscopes are often used to stop internal bleeding by clamping together the edges of wounds or incisions.\nIn the simplest form, these clips grasp tissue surrounding a wound, bringing edges of the wound together to allow natural healing processes to close the wound. Specialized endoscopic clipping devices are used to deliver the clips to desired locations within the body and to position and deploy the clips at the desired locations after which the clip delivery device is withdrawn, leaving the clip within the body.\nEndoscopic hemostatic clipping devices are generally designed to reach tissues deep within the body (e.g., within the GI tract, the pulmonary system, the vascular system or other lumens and ducts) via a working lumen of an endoscope. Thus, the dimensions of the clipping device are limited by the dimensions of the working channels of endoscopes with which they are to be employed."}
{"text": "1. Field of the Invention\nThe present invention relates to drainage for land areas. More particularly, the present invention relates to a system for draining land areas, such as golf courses and farms through the collection of water into a principal catch basin having permeable walls, and subsequently siphoning the water collected from the catch basin to a distant exit point such as a lake, canal, ditch or the like. There could be further included a series of secondary catch basins draining into the principal permeable catch basin, so that an even larger area may be drained through siphoning from a single control system.\n2. General Background\nIn the present state of the art, drainage water such as water from heavy rains or the like may be collected either by a catch basin or a sub-surface drainage system that collects water through seepage. At the present time, a catch basin, which consist of a substantially box-like chamber is set within the ground and positioned at a low point where water would flow under gravity and be collected within the catch basin. At this point the water is exited through a solid pipe at a one percent slope whenever possible, so that the water then exits into some type of collection system such as a lake, pond or creek. However, often times there is insufficient fall from the point of the collection system to the exit point to either obtain a 1% grade or any positive grade at all. The only alternatives, therefore, in the past, have been to either not attempt to do the drainage at all or to attempt to build a gravel sump, dry well, or sump pump. The concept behind these systems relate to the fact that the gravel sump will provide some temporary storage capacity for the water moving out of the collection system and maintain it until it can slowly seep back into the soil profile. Often times there will be a dry well built that will extend deep enough to reach a permeable strata below the less permeable layers that are near the surface. The system has two main problems, i.e., (a) the sumps have very little capacity, and once the storage capacity is filled, the collection system will back up and stop functioning; and (b) for a system to extend deep enough into the ground to hit a sand strata, it can be a very expensive system to install, depending on the depth of the sand strata. The Environmental Protection Agency is now concerned that bringing surface water down to these strata may contaminate the water table.\nThe other alternative is to build a sump pump which operates a pump off of a float switch and pumps the water out every time the level of the water activates the float switch. The problem with these systems is that they require a large amount of maintenance, are expensive, and a pump has to be located at each sump.\nThere also is a problem with the existing type of catch basin that is presently installed. These systems operate very well when they are surrounded by an impermeable surface, such as asphalt or concrete for street or parking lot drainage. However, in land drainage such as farms and golf courses, these systems are surrounded by permeable areas of soil. The problem is that before the water can enter the catch basin system, many times a large amount of water is trapped into the permeable soil profile surrounding the catch basin itself. Once that water is trapped into the soil profile, it cannot enter through the solid walls constructed of concrete or brick, and stays in the soil profile, thus creating a wet soggy saturated area. The other problem with existing catch basins is that they are designed only to stop the largest particles from entering the pipe that is used to drain the water out of the catch basin. Therefore, many contaminants such as small rocks and branches become lodged into the piping system. To prolong the time period before clogging occurs in drainage pipes, it has been a common practice to oversize the pipe in relation to the amount of water that is necessary to be drained. This oversizing of pipe to prevent clogging increases the expense of materials for the piping, as well as the installation costs because larger amounts of dirt must be excavated to install the pipe. Also, there is no system provided in existing drainage systems for allowing the escape of air trapped within the soil profile. This air reaches an equilibrium with the infiltrating water and stops the infiltration of the water into the soil profile. If this air were allowed to escape, a larger portion of the profile could absorb the infiltrating rainwater.\nThere have been several patents which address the drainage of water in a system. The most pertinent are listed as follows, including a narrative of each submitted in applicants Prior Art Statement, filed herewith:\n______________________________________ U.S. Pat. No. INVENTOR TITLE ______________________________________ 4,299,697 Curati, Jr. \"A Liquid Containment And Storage System For Rail- road Tracks\" 1,734,777 Pike \"System Of Draining\" 462,864 Hershberger \"Catch Basin For Sewers\" 2,432,203 Miller \"Catch Basin\" 4,472,086 Leach \"Geotextile Fabric Construction\" ______________________________________"}
{"text": "Brief Description of the Prior Art\nResilient contacts for connecting a conductor to a bus bar are well known in the patented prior art, as evidenced by the U.S. patents to Beege, et al., U.S. Pat. No. 5,938,484 and Delarue, et al., U.S. Pat. No. 5,879,204, among others.\nAs evidenced by the German Auslegeschrift No. 1,261,923, the German patent No. 4233446 and the European patent No. 908,965, the clamping edge at the free end of the clamping leg of the resilient contact normally continues to be pointed toward the electrical conductor as the clamping leg is displaced toward the released position, whereupon as the conductor is withdrawn from the contact, there is a possibility of damaging the conductor, which may be particularly undesirable in the case of a multi-wire conductor, or a fine-wire conductor. This applies also to the resilient contact of the German patent No. 2062158 B2, wherein the free end of the clamping leg of the resilient contact is bent only at a small angle.\nAs shown by the German Auslegeschrift No. 1,261,923 and the European patent No. 908,8965, it has been proposed to provide adjacent the clamping leg of a resilient contact an operating member that functions to displace the clamping leg toward a disengaged position relative to the conductor. These devises have generally proven to be undesirable, however, owing to the fact that the actuation surface on the clamping leg lies ahead of the clamping edge. When the actuating devices are in the form of levers or rams, as shown by the German patent Nos. DE 2062158 B2 and DE 4233446 C1, for example, the points of engagement are at undesirable locations on the clamping leg ahead of the clamping edge.\nThe present invention was developed to avoid the above and other drawbacks of the known resilient connectors."}
{"text": "The treatment of bacterial infections with antibiotics is one of the mainstays of human medicine. Unfortunately the effectiveness of antibiotics has become limited due to an increase in bacterial antibiotic resistance in the face of a decreasing efforts and success in discovery of new classes of antibiotics. Today, infectious diseases are the second leading cause of death worldwide and the largest cause of premature deaths and loss of work productivity in industrialized countries. Nosocomial bacterial infections that are resistant to therapy result in annual costs of more than $2 billion and account for more than 80,000 direct and indirect deaths in North America alone, whereas a major complication of microbial diseases, namely sepsis, accounts for 700,000 cases and 140,000 deaths in North America.\nA major limitation in antibiotic development has been difficulties in finding new structures with equivalent properties to the conventional antibiotics, namely low toxicity for the host and a broad spectrum of action against bacterial pathogens. Recent novel antibiotic classes, including the oxazolidinones (linezolid), the streptogramins (synercid) and the glycolipids (daptomycin) are all active only against Gram positive pathogens. Cationic antimicrobial peptides, found in most species of life, represent a good template for a new generation of antimicrobials. They kill both Gram negative and Gram positive microorganisms rapidly and directly, do not easily select mutants, work against common clinically-resistant bacteria such as methicillin-resistant Staphylococcus aureus (MRSA) and vancomycin resistant Enterococcus (VRE), show a synergistic effect with conventional antibiotics, and can often activate host innate immunity without displaying immunogenicity (Hancock R E W. 2001. Cationic peptides: effectors in innate immunity and novel antimicrobials. Lancet Infectious Diseases 1, 156-164). Moreover, they seem to counteract some of the more harmful aspects of inflammation (e.g. sepsis, endotoxaemia), which is extremely important since rapid killing of bacteria and subsequent liberation of bacterial components such as LPS or peptidoglycan can induce fatal immune dysregulation (Jarisch-Herxheimer reaction) (Gough M, Hancock R E W, Kelly N M. 1996. Anti-endotoxic potential of cationic peptide antimicrobials. Infect. Immun. 64, 4922-4927). A need exists in the art for developing new treatments for infections to be used as broad spectrum antibiotics and/or as agents that selectively enhance aspects of innate immunity while suppressing potentially harmful inflammation.\nThe innate immune system is a highly effective and evolved general defense system that involves a variety of effector functions including phagocytic cells, complement, and the like, but is generally incompletely understood. Elements of innate immunity are always present at low levels and are activated very rapidly when stimulated by pathogens, acting to prevent these pathogens from causing disease. Generally speaking many known innate immune responses are “triggered” by the binding of microbial signaling molecules, like lipopolysaccharide (LPS), with pattern recognition receptors such as Toll-like receptors (TLR) on the surface of host cells. Many of the effector functions of innate immunity are grouped together in the inflammatory response. However, too severe an inflammatory response can result in responses that are harmful to the body, and, in an extreme case, sepsis and potentially death can occur; indeed sepsis occurs in approximately 780,000 patients in North America annually with 140,000 deaths. Thus, a therapeutic intervention to boost innate immunity, which is based on stimulation of TLR signaling (for example using a TLR agonist), has the potential disadvantage that it could stimulate a potentially harmful inflammatory response and/or exacerbate the natural inflammatory response to infection. A further need exists in the art for therapeutic inverterventions to boost innate immunity that are effective and have fewer undesirable side effects or adverse reactions."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of quantitative analysis of hexavalent chromium in a coating for a metal substrate such as a chromate coating used in household electronic appliances and automobiles. The present invention also relates to a method for controlling a hazardous element in an encapsulating resin of a resin encapsulation semiconductor device with a fluorescent X-ray analyzer.\n2. Description of the Related Art\nIn Europe, it is required to restrict the use of lead, mercury, cadmium, PBB, PBDE, and hexavalent chromium in principle in accordance with the Restriction of Hazardous Substances Directive (ROHS) which takes effect on Jul. 1, 2006. In order to comply with the directive, it is desired to develop analytical methods capable of conveniently assaying these substances.\nIt is known that X-ray photoelectron spectroscopy can be used as an analytical method for directly assaying hexavalent chromium content in a chromate coating (see Jpn. Pat. Appln. KOKAI No. H05-164710). Since this method assays only the surface of the coating, however, there are problems that the analytical value does not represent the value for the entire sample and that it is hard to separate peaks of trivalent chromium and hexavalent chromium.\nISO 3613 defines an analytical method for determining hexavalent chromium eluted from a chromate coating using boiling water. This method assays only hexavalent chromium eluted by the boiling water, and cannot determine total hexavalent chromium content present in the chromate coating.\nAn analytical method using sodium hydroxide solution as an extracting solution is known in “Method 3060” defined by the United States Environmental Protection Agency (EPA). When a chromate coating on a substrate of an amphoteric metal such as aluminum is assayed using the extracting solution, the aluminum substrate is eluted prior to the chromate coating accompanied by hydrogen gas generation and reduction of hexavalent chromium to trivalent chromium, making it difficult to perform highly accurate quantitative analysis of hexavalent chromium.\nFurther, as an analytical method of efficiently extracting hexavalent chromium in a chromate coating on a metal substrate in a short time, a method is known in which the chromate coating is cracked and then the coating is immersed in an extracting solution to extract hexavalent chromium so as to be analyzed (see Jpn. Pat. Appln. KOKAI No. 2004-325321). Since this method requires to applying thermal shock or mechanical shock to the substrate in order to cause cracks, however, there is a possibility that the substrate itself may be broken. Thus, this method cannot be generally used for analyzing hexavalent chromium.\nIn addition, according to the RoHS directive, it is required to determine that particular hazardous substances are not contained in an electronic material. In compliance with the RoHS directive, it is expected to develop an analytical method capable of measuring hazardous elements such as Br, Sb, As, Bi and Pb easily and precisely and a method for controlling the hazardous elements precisely using the analytical method.\nIn order to non-destructively analyze hazardous elements in an encapsulating resin of a resin encapsulation semiconductor device, use of a fluorescent X-ray analyzer is easy and effective. The fluorescent X-ray analysis generally employs a fundamental parameter (FP) method for a metal sample and a calibration curve method for a resin sample. In order to prepare the calibration curve, standard resin samples are used to which a known quantity of element to be analyzed is added. The standard resin samples having a matrix resin such as polyethylene, ABS resin and vinyl chloride resin are commercially available.\nSince the encapsulating resin is compressed in molding, it has a higher density than the standard resin samples. In addition, the encapsulating resin contains filler such as silica for enhancing strength. Thus, a primary X-ray is hard to penetrate into the encapsulating resin, which tends to bring about a lower analytical value of a fluorescent X-ray. Further, when the encapsulating resin of the resin encapsulation semiconductor device is analyzed with a fluorescent X-ray analyzer, there is a possibility to detect a fluorescent X-ray of constituent elements of a lead frame, semiconductor chip and wire, which may be a cause of an error for a fluorescent X-ray of some hazardous element to be detected.\nConventionally, there is proposed a method for judging whether a sample such as polystyrene contains lead (Pb) by fluorescent X-ray analysis. See JP-A 2007-3331 (KOKAI). The method comprises applying an X-ray to a sample, preparing a fluorescent X-ray spectrum, and judging that the sample contains Pb when the spectrum has peaks at all energy positions corresponding to Pb. The method aims at analyzing Pb precisely by avoiding influence of As or Br having a peak overlapping the peak of Pb. However, the method cannot solve the problem of precision due to the lower analytical value for an encapsulating resin of a resin encapsulation semiconductor device which is compressed and has a high density as well as contains filler.\nAlso, there is known a method for judging a material containing a particular substance in an object to be measured in which materials are formed in layers by fluorescent X-ray analysis. See JP-A 2007-17306 (KOKAI). This publication discloses a sample of Fe on which Zn plating and Cr plating are deposited as the object of layered structure and particular substances such as Cd, Pb, Hg, Br, Cr, Au, Ag, Pt and Pd. The method comprising applying an X-ray to a sample at a controlled depth by varying X-ray irradiation conditions, detecting a fluorescent X-ray emitted from the sample, and judging the material in the layered structure in which the particular substance is contained based on synchronous detention of the material and particular substance. The method is effective for a sample having a relatively simple layered structure and capable of providing sufficient fluorescent X-ray intensity. However, the method cannot solve the problem of precision due to the lower analytical value for an encapsulating resin of a resin encapsulation semiconductor device which is compressed and has a high density as well as contains filler. In addition, the method is not suited for analyzing an encapsulating resin of a resin encapsulation semiconductor device in which a lead frame, chip and wire are arranged into a complicated structure in a matrix resin."}
{"text": "Recent advances in hardware and communication technologies have brought about an era of client computing over wired and wireless networks. The proliferation of powerful notebook computers and wireless client devices promises to provide client end users with network access at any time and in any location over various networks, including the Internet. This continuous connectivity allows users to be quickly notified of changing events, and provides them with the resources necessary to respond in realtime even when in transit. For example, in the financial services industry, online traders and financial professionals may be given the power to access information in real-time, using wireless client devices.\nConventionally, however, developers of complex, wireless messaging solutions have been forced to develop applications that are specific to various device types and network access protocols in diverse enterprise architectures and platforms. In other words, conventional client computing solutions have been largely platform-specific, network-specific, or both. For example, messages may be generated by a wide variety of applications running on a wide variety of client devices, such as Palm computing platform devices, Windows CE devices, paging and messaging devices, laptop computers, data-capable smart phones, etc. Depending on the type of network used by service providers, the client-generated messages may be transported over networks having various access protocols, such as, e.g., Cellular Digital Packet Data (CDPD), Mobitex, dial-up Internet connections, Code Division Multiple Access (CDMA), Global System for Mobile Communications (GSM), and ReFlex. As a result, current developers of client computing solutions must have intimate knowledge of specific network characteristics including, e.g., wireless network characteristics, protocol environments, and wireless links channel characteristics. Therefore, there exists a need to simplify wireless client and server application development environments to support the wide variety of device and network dependent architectures.\nMessaging Application Programming Interface (MAPI) is a messaging architecture and an interface component for applications such as electronic mail, scheduling, calendaring and document management. As a messaging architecture, MAPI provides a consistent interface for multiple application programs to interact with multiple messaging systems across a variety of hardware platforms. MAPI provides cross platform support through such industry standards as Simple Mail Transfer Protocol (SMTP), X.400 and common messaging calls. MAPI is also the messaging component of Windows Open Services Architecture (WOSA).\nAccordingly, MAPI is built into such operating systems as, e.g., Windows 95, Windows 98, Windows NT and Windows 2000, available from Microsoft Corporation of Redmond, Wash., U.S.A. and can be used by 16-bit and 32-bit Windows applications. For example, a word processor can send documents and a workgroup application can share and store different types of data using MAPI. MAPI separates the programming interfaces used by the client applications and the service providers. Every component works with a common, Microsoft Windows-based user interface. For example, a single messaging client application can be used to receive messages from fax, a bulletin board system, a host-based messaging system and a LAN-based system. Messages from all of these systems can be delivered to a single “universal inbox.”\nTransmission Control Protocol (TCP) is a transport layer protocol used by an application in one host to communicate with an application in another host. This is accomplished by services provided by the protocol layers beneath the transport layer in both hosts. As a connection-oriented protocol TCP requires the establishment of a connection between the two hosts before two applications are able to communicate. TCP manages the connection and once both applications have communicated all required information between themselves the connection is released as if the connection is two simplex connections as opposed to a single duplex connection. The information is transferred between applications on different hosts is a byte stream. The transport layer hides message transfer details such as segmentation, ordering and duplication from the applications and provides end-to-end acknowledgement.\nThe Internet Protocol (IP) layer provides certain services to the transport layer protocol including hiding the details of the physical and data link layers and the services provided by them from the transport layer protocol. The IP layer provides a datagram delivery service. A datagram is a unit of data less than an entire message. A message may be, for example, a file, which may be quite large. Since there is a maximum size for a message (or file), the message may have to be segmented and transferred in smaller units. These smaller units are thus called datagrams. Each datagram is sent over the network as a separate entity and may, in fact, follow separate paths to the destination host. At the destination host, the datagrams are reassembled in proper order (usually in a buffer area) by the transport layer. Each node on the network sends any datagrams on to a next node only considering the final destination and only acknowledges receipt of the datagram to the preceding node. That is, the IP layer does not provide end-to-end acknowledgement. End-to-end acknowledgement is a service of the transport layer protocol. Should any node-to-node acknowledgements not be received by the preceding node, the datagram or datagrams unacknowledged will be retransmitted. The transport layer in the destination host will also acknowledge any duplicated datagrams (else receipt of duplicate datagrams will continue resulting in a clogged network) and ignore them.\nRouting between network nodes is accomplished by means of routing tables. These routing tables can be static or dynamic and result in datagrams being forwarded from a source host to a destination host one node at a time. The intermediate nodes are often called “hops.”\nThe acronym, TCP/IP, is also used to refer to a five layer protocol model similar to the ISO/OSI seven layer protocol model. The TCP/IP model does not have the equivalent to layers 5 and 6 of the ISO/OSI protocol model. A protocol model defines the protocol layers and the interfaces between the layers. When implemented in software, hardware or firmware or possibly field programmable gate arrays (FPGAs), the implementation provides the actual services. This layered approach allows for ease of upgrading so long as the interface to the layer immediately above or below is not altered. Layering also allows for complete substitution. For example, should a new physical medium become available then so long as the interface between layer two and layer one remain the same, an old physical layer implementation module can be removed and a new implementation module substituted. In the alternative, the new implementation module could be added as another option. That is, the protocol suite defines the rules and the implementation provides the services that allow the communications to take place between applications on different hosts. The implementation of the TCP layer provides for the application to require a certain Quality of Service (QOS) as specified by a set of parameters including but not limited to priority, transmission delay, throughput etc.\nAnother well-known transport layer protocol is known as User Datagram Protocol (UDP), which is a connectionless transport protocol. The basic data transfer unit is the datagram. A datagram is divided into header and data areas as it is for the IP layer. An additional header over and above the IP header is used. The UDP header contains source and destination addresses (ports), a UDP length field (the length includes the 8 byte header and the data) and a UDP checksum. The UDP data includes the IP header and data. The IP layer supports UDP as a connectionless transport protocol for use in transmitting, for example, one time request-reply type messages or applications where time is of greater importance than accuracy.\nTCP is used by applications on different hosts to communicate over an assumed unreliable network. To accomplish such communication much is added to the protocol in order to ensure that the information transferred over the network is valid. This addition has a cost and that cost is increased overhead with the attendant increase in bandwidth. A UDP header is eight bytes, the TCP header is 24 bytes and an IP header is a minimum of 20 bytes. Therefore, UDP/IP headers are a minimum of 28 bytes and TCP/IP headers are a minimum of 44 bytes. This is fairly large in terms of overhead and bandwidth utilization particularly over wireless networks. There are other significant problems with using standard TCP/IP over wireless networks principally in the area of flow control. The UDP/IP protocol combination, while not offering any guarantees to users, is expected to be reliable. Wireless networks tend, however, by their nature to be lossy. Several solutions have been proposed when the network is not homogeneous. That is, when the network has both wireless and wireline portions. One suggestion is to use indirect TCP and another is to use snooping.\nOther protocols such as Serial Line IP (SLIP) and Point to Point Protocol (PPP) have been developed. SLIP is not a standard and both are for point to point connections only so are not available for uses in networks. CDPD vendors indicate that they provide an integrated TCP/IP stack but it is not known the cost in terms of bandwidth overhead.\nConventionally, the existing wireless protocols do not provide an end-to-end solution over multiple networks and multiple client devices. Therefore, in addition to the need for a common architecture through a single, user friendly methodology for providing effective and reliable wireless data solutions for hand-held and laptop computing devices, wireless networks, and legacy systems, there also exists a need to efficiently and reliably communicate data using minimum bandwidth."}
{"text": "The present invention relates to a display control device and a data processing unit using the same, and a technique useful in application to a communication portable terminal device having a liquid crystal controller driver.\nA liquid crystal controller driver performs control for writing display data supplied from a host device into a frame buffer memory out of synchronization with its internal operation timing, for reading display data written into the frame buffer memory in synchronization with the internal operation timing, and for driving signal electrodes of a liquid crystal panel while using the read display data to synchronize with the timing of display. In Japanese Unexamined Patent Publication No. JP-A-2003-131622, there is description about a technique for display data compression and storing on a frame buffer memory with the aim of increasing the CPU throughput required for write and read of display data on the frame buffer memory. Specifically, iterations of image data in display data are counted in units of pixels, and the top of the iterations and the number thereof are stored in combination as information of iteration. Then, using the information of iteration in reading display data, the readout of image data iterated can be omitted."}
{"text": "An apparatus designed to obtain mechanical energy from the atmosphere or water.\n1. Field of the Invention\nThe Atmospheric/Aqua Turbine is most closely related to the Wind-Charger, the Windmill and the Hydroelectric Water Turbine.\nThe Atmospheric/Aqua Turbine is an apparatus that can be incorporated into a vehicle for the purpose of extracting mechanical energy from the resistance of the atmosphere or water when the vehicle or boat it is installed on is in motion.\n2. Description of Prior Art\nThe Wind-Charger, the Windmill and the Hydroelectric Water Turbine all three convert wind or water velocity into mechanical energy at fixed locations. They all depend on the wind blowing or water flowing in order to produce mechanical energy which is used to produce electricity or to pump water.\nThe Atmospheric/Aqua Turbine is designed to extract mechanical energy from the atmosphere or water by being incorporated into the configuration of a vehicle. The Atmospheric/Aqua Turbine has a scoop which scoops up atmospheric pressure or water and passes it through two fan chambers in order to extract mechanical energy from the atmosphere or water. The Atmospheric/Aqua Turbine will operate when the vehicle it is installed on is moving and depends on that movement instead of depending on the wind to blow or water to flow."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus or the like which is used in a printer, a copier, a facsimile or the like, and more particularly to an image forming apparatus or the like in which plural image carriers are arranged in parallel.\n2. Description of the Related Art\nAs an image forming apparatus using an electrophotographic method, a so-called tandem-type image forming apparatus in which plural image forming units of yellow, magenta, cyan and black are arranged in parallel is popularly used. In the tandem-type image forming apparatus, in general, a photosensitive drum having a photosensitive layer on an outer peripheral surface thereof is provided to each image forming unit. At a position which faces the photosensitive drums in an opposed manner, a transfer belt which is extended over plural expansion rolls is arranged. Toner images of respective colors which are formed over the respective photosensitive drums are sequentially transferred onto the transfer belt directly or by way of a recording sheet (various types of sheets including a sheet, an OHP sheet or the like) in a multiplied manner. Further, there also exists an image forming apparatus in which a transfer roll which is brought into contact with each photosensitive drum is provided for every image forming unit and a transfer action of a recording material is performed by the photosensitive drum and the transfer roll.\nIn this tandem-type image forming apparatus, it is necessary to transfer respective color-component images formed by the respective image forming units to a same image carrier position of the transfer belt or a same recording material in a sequentially overlapped manner and hence, there arises a drawback that the registration of the respective color-component images is difficult. Accordingly, with respect to techniques described in publications on the related art, there exists a technique in which a single drive motor is used as a drive system and plural photosensitive drums which are arranged in parallel are driven by the single drive motor. To be more specific, there exists the technique in which plural worm gears are mounted on a rotary drive shaft which is connected to and rotatably driven by the drive motor, worm wheels are mounted on rotary shafts of the respective photosensitive drums, and the worm gears and the worm wheels are meshed with each other whereby the plural photosensitive drums are driven by the single drive motor (see, for example, JP-A-63-11967 (pages 3 to 4, FIG. 1), also referred to as “patent literature 1” hereinafter).\nAccording to the technique described in the patent literature 1, compared to a case in which the respective photosensitive drums have drive sources individually, the registration of the respective color-component images is facilitated. However, mainly attributed to meshing having a backlash between the worm gear and the worm wheel, there has been found a technical drawback that the speeds of the photosensitive drums vary. In view of the above, there has been proposed, for example, a technique which provides a centering position constraining device for rotatably holding one ends of the worm shafts while centering one ends of the worm shafts and constraining both end positions of each worm shaft with respect to the axial direction whereby the image irregularities attributed to the speed variation of the photosensitive drums can be suppressed (see, for example, JP-A-2000-131918 (pages 4 to 5, FIG. 1), also referred to as “patent literature 2” hereinafter).\nIn this manner, by adopting the technique described in JP-A-2000-131918, it is possible to suppress the speed variation attributed to the oscillation variation and the axial oscillation of the worm shafts. However, the related art described in the patent literature 2 is not provided to suppress the speed variation caused by geometric factors such as the tooth shape accuracy, the pitch accuracy and the like of the worm gear as a single body. To maintain a certain rotation with high accuracy by suppressing the speed variation attributed to the geometric factors such as the tooth shape accuracy, the pitch accuracy or the like, it is necessary to manufacture the worm gears per se with high accuracy.\nIn manufacturing the general worm gear used in the tandem-type image forming apparatus, a metal material such as stainless steel is formed into a worm shape by cutting machining, rolling forming or the like and, thereafter, for ensuring the accuracy of the gear, a gear grinding step for grinding teeth of the gear is provided separately. Accordingly, to obtain the tooth shape having high accuracy, the manufacturing cost is eventually pushed up whereby there arises a drawback that a product cost is extremely pushed up."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method for determining variation in a predetermined physical property of a circuit.\n2. Description of the Prior Art\nIn modern data processing systems, the individual components of data processing circuits are becoming smaller and smaller. As process geometries shrink, the variations in certain physical properties of the individual circuit components is becoming more significant, and techniques are required to enable those variations to be determined.\nOne physical property that may vary significantly between circuit components is referred to as an ageing effect, a variation in the ageing effect amongst components resulting from a variation in performance degradation of individual transistors over time. Examples of articles that describe monitoring techniques for the ageing effect are “An Array-Based Odometer System for Statistically Significant Circuit Aging Characterization” by John Keane et al, IEEE Journal of Solid-State Circuits, Volume 46, No. 10, October 2011, and the article “A Precise-Tracking NBTI-Degradation Monitor Independent of NBTI Recovery Effect” by E Saneyoshi et al, 2010 IEEE International Solid-State Circuits Conference, Session 9, Digital Circuits and Sensors. However, both of these techniques are invasive, requiring an external tester to be used.\nAnother type of physical property that may vary between components within a circuit is a history effect that can occur in transistors that have a body region insulated from the substrate, the body region comprising the channel material in which a channel is established between a source and drain of the transistor below the gate of the transistor. One known technology that can be used to produce transistors having a body region insulated from the substrate is referred to as Silicon-On-Insulator (SOI) technology, where the SOI channel material is formed within a region of a thin superficial silicon layer above an oxide insulating layer and located under the gate of the transistor, reducing resistive leakage to the substrate and capacitive loading by the substrate. However, it has been found that the voltage on the body region becomes dependent on the previous circuit activity, this typically being referred to as the history effect. This variation in voltage on the body region can have an impact on the behaviour of components constructed using such transistors. The article “A New Structure for In-Depth History Effect Characterization on Partially Depleted SOI Transistors” by O Faynot et al, 2002 IEEE International SOI Conference, 10/02, pages 35 to 36, describes a technique for measuring the history effect, but similarly to the earlier described articles is invasive, and requires use of an external test probe.\nWhilst techniques that use external testing equipment can be useful during the design and manufacture of components, for example to allow suitable margins to be set, it is still desirable to provide mechanisms that allow variations in certain physical properties to be monitored “on the fly” whilst circuits are in use, since if accurate information about such variations can be obtained in real time during use of the circuit, it may be possible to take certain actions in order to reduce or at least constrain such variations. The article “Run-Time Adaptive Performance Compensation using On-Chip Sensors” by M Hashimoto, Department of Information Systems Engineering, Osaka University and JST, CREST, ASPDAC '11 Proceedings of the 16th Asia and South Pacific Design Automation Conference, pages 285-290, describes a run-time performance adaptation technique using on-chip sensors for capturing delay variations, with circuit performance then being compensated according to the sensor outputs. The described technique does not provide any quantitive information as to the variations being observed.\nA number of articles have been published that discuss techniques for monitoring the ageing effect of transistors on the fly during operation of integrated circuits. Examples of such articles include the article “A Compact Monitoring Circuit for Real-Time On-Chip Diagnosis of Hot-Carrier Induced Degradation” by H Oner et al, Proceedings of the IEEE 1997 International Conference on Microelectronic Test Structures, Volume 10, March 1997, the article “A TDC-Based Test Platform for Dynamic Circuit Ageing Characterisation” by M Chen et al, 2011 IEEE International Reliability Physics Symposium (IRPS), 10-14 Apr. 2011, 2B.2.1-2B.2.5, the article “On-Chip Circuit For Monitoring Frequency Degradation due to NBTI” by K Stawiasz et al, IEEE CFP08RPS-CDR 46th Annual International Reliability Physics Symposium, Phoenix, 2008, Pages 532 to 535, and the article “An All-In-One Silicon Odometer for Separately Monitoring HCI, BTI and TDBD” by J Keane et al, IEEE Journal of Solid-State Circuits, Volume 45, No. 4, April 2010, Pages 817 to 829. Each of these techniques uses a ring oscillator as the circuitry to be monitored, and as such lacks flexibility in that the circuitry being monitored is somewhat artificial and cannot directly represent an actual circuit portion of a real data processing apparatus, for example a critical path within a processor. In addition, such techniques using ring oscillators as the monitored circuit cannot measure certain physical properties such as the earlier-mentioned history effect, since the history effect occurs primarily due to static data and hence the continual dynamic changes that occur within an oscillator will mask the history effect.\nIn the article “A Test Concept For Circuit Level Aging Demonstrated By A Differential Amplifier” by F. Chouard et al, 2010 IEEE International Reliability Physics Symposium, (IRPS), 2-6 May 2010, pages 826-829, the authors propose applying an accelerated aging technique to an analogue circuit, thereby allowing measurements to be taken from a tester after forced aging of the circuit to model the aging effects. This approach hence seeks to take aging into account by modeling/simulating aging rather than by measuring actual aging in a live circuit.\nAccordingly, it would be desirable to provide a more flexible approach for determining variation in a predetermined physical property of a circuit, that can be used on the fly as and when required during the lifetime of the circuit."}
{"text": "It is known that, since the start of civil aviation, the take-off of an airplane is handled manually by the airplane's pilot by means of piloting units (central control column or side-stick controller, rudder bar) that are specific to the airplane concerned. During this take-off phase, the crew is responsible for controlling the path of the airplane, both in the vertical plane and in the lateral plane (or horizontal plane). The crew has to follow theoretical guidance objectives (recommendation of the airplane's flight manual), visual guidance objectives (axis of the runway) and/or physical guidance objectives (for example, reference of an LOC-type antenna, specified below).\nIn such a standard manual take-off, there arises the problem of the quality with which the guidance objectives are followed during the take-off.\nIf we look firstly at the elevation guidance (that is, in the vertical plane), the fact that the crew is responsible for the take-off renders the quality of the rotation directly dependent on the behavior of the pilot.\nAssessing the rotation speed, the speed at which the pilot must pull on the control column (central control column or side-stick controller) to initiate the rotation of the airplane, is done visually.\nBecause of this, a greater or lesser time delay may occur between the moment when the pilot realizes that he has reached the rotation speed and the moment when he pulls on the control column.\nAt this instant in the take-off, the airplane is at high speed, and this time delay is reflected in a dispersion over the take-off distance.\nIn the opposite case, it may be that the pilot pulls too early on the control column. In this case, the speed of the airplane does not allow it to take off, and the airplane continues to accelerate with the nose lifted. This is reflected in a very significant drag, and also, a direct impact on the take-off distance.\nFurthermore, with large airplanes, a risk of tailstrike on take-off must be taken into account. It is, in fact, not rare to observe this type of incident in service, which normally results in a lengthy lay-up of the airplane. Such an incident is normally the result of a nonconforming action on the part of the pilot (excessive pull-up command) or an unforeseen external phenomenon (violent gust of wind, for example).\nIf we now consider, generally, the lateral guidance of the airplane, there is today just one source of help in the lateral guidance, namely a lateral alignment beam of LOC type which is emitted from the ground and which is detected on the airplane using a detection system. Such a lateral alignment beam or LOC beam is normally emitted by a directional VHF radiotransmitter, which is placed on the axis of the runway, at the end opposite from the threshold. This LOC beam is intended for landing and, in principle, ensures azimuth guidance along the approach axis, and this according to an ideal lateral alignment profile in a precision instrument approach of the ILS (Instrument Landing System) type. This radiotransmitter normally emits two signals with different modulation which overlap in the approach axis where the two-signals are received with equal intensity. Although this LOC beam can also be used for the take-off, it was created, and is characteristics have therefore been more particularly adapted, for use on a landing.\nIf a sight guidance mode (as mentioned above) remains acceptable in clear weather, it becomes very difficult in reduced visibility conditions. Similarly, the workload of the pilot greatly increases if an engine fails. The pilot must then trim the thrust difference while ensuring a good rotation.\nThus, the performance levels of the airplane on take-off directly depend on the behavior and the feelings of the pilot (stress, fatigue). This phenomenon is amplified when the weather conditions worsen. The wind is in particular the most disturbing factor in this phase."}
{"text": "Field\nThe present disclosure generally relates to computing devices, and more particularly to the provision of data on a computing device.\nDescription of the Related Art\nComputing systems commonly load (e.g., download) software updates in order to keep an installed application current and secure. In many instances, the update or other modification may be loaded for an operating system of a computer. Certain systems permit the update or other modification to be uninstalled or otherwise unloaded if the update or other modification is insubstantial in order to return the system to a previous state. The process and time required for unloading such an insubstantial update or modification is nonetheless commonly very long, and in some instances the unloading may be unsuccessful. In instances where the update for the system is substantial, many systems do not permit the update to be uninstalled or otherwise unloaded because the system has changed too extensively as a result of the update such that the system cannot return to a previous state."}
{"text": "1. Field of the Invention\nThe present invention generally relates to the field of effect devices that can connect with other effect devices in parallel.\n2. Related Art\nVarious types of effect devices (effectors) that add effects to input musical sound signals and then output the resulting sounds, have been known. Such effect devices include compact effectors which add different effects for the output of electric guitars equipped with pedal switch, or the multiple effector which can use a combination of multiple effects.\nFIG. (Diagram) 7 is a block diagram that depicts an electrical configuration of a traditional Effect Device 30. This device has an input Jack JK 2; two amplifiers, AMP 1 and AMP 2; one Switch SW 1; an effect add circuit 2; and an output Jack JK 1. A plug on one end of a connecting cord for an electrical guitar, synthesizers, or the like may be inserted in the JK2 input Jack. Musical sound signals that are generated by sources such as an electric guitar or a synthesizer, are used as the input.\nThe musical sound signals that are input through the input Jack (JK 2), are provided to the effect add circuit 2, via the AMP 2 amplifier, and are connected to one side of the input terminal (1 Side of SW1 switch). The output of the effect add circuit 2 is connected to the other side of the input terminal (2 Side of Switch SW 1). The output terminal of Switch SW 1 may be connected to the output Jack JK 1, via AMP-1 Amplifier. The effect add circuit 2 can add several types of effects including a distortion group effect which distorts the music sound, and a resonance effect which reverberates or delays the musical sounds.\nA plug on one end of a connecting cord that connects with the operating amplifier of a speaker (not shown on Diagrams) is inserted to output Jack JK 1. The SW 1 switch is equipped with foot pedals that a player can operate with his/her foot. The player can operate the foot pedals with his/her foot while he/she is playing guitar or synthesizer with both of his/her hands.\nWhen Switch SW 1 is operated to connect with 1 Side, the input musical sound signals will be output from the output Jack (JK 1) without any effect. When the Switch SW1 is connected to Side-2, the effect add circuit 2 effects the input musical sound signals and the effected musical sound signals are output from the output Jack (JK 1).\nTraditional effect devices such as the one explained above, are equipped with only one input jack and one output jack respectively. Such devices may be connected together in series, where an output of primary effect device is input on subsequent effect device. However, there can be difficulties in trying to achieve a parallel connection of such devices where musical sound signals are input to multiple effect devices and the output of the effect devices are mixed without one effect device becoming dominant over the other(s). A parallel connection implementation needed a means to separate musical sound signals on the input side of the effects devices, and a means to mix musical sound signals on the output side of the effects devices. The configuration of such a system has been known to be relatively complex. Serial connection has been an industry preference for this reason. When such devices are connected in series, the effect from a subsequent effect device tends to become dominant. For this reason it can be impossible to have certain desired effects with the serial connections of such prior devices."}
{"text": "Wireless access systems have been widely deployed to provide various types of communication services such as voice or data. In general, a wireless access system is a multiple access system that supports communication of multiple users by sharing available system resources (a bandwidth, transmission power, etc.) among them. For example, multiple access systems include a Code Division Multiple Access (CDMA) system, a Frequency Division Multiple Access (FDMA) system, a Time Division Multiple Access (TDMA) system, an Orthogonal Frequency Division Multiple Access (OFDMA) system, and a Single Carrier Frequency Division Multiple Access (SC-FDMA) system.\nAs a number of communication devices have required higher communication capacity, the necessity of the mobile broadband communication much improved than the existing radio access technology (RAT) has increased. In addition, massive machine type communications (MTC) capable of providing various services at anytime and anywhere by connecting a number of devices or things to each other has been considered in the next generation communication system. Moreover, a communication system design capable of supporting services/UEs sensitive to reliability and latency has been discussed.\nAs described above, the introduction of the next generation RAT considering the enhanced mobile broadband communication, massive MTC, Ultra-reliable and low latency communication (URLLC), and the like has been discussed."}
{"text": "Credit cards, debit cards, prepaid transaction cards, gift cards, and other similar electronic tokens are convenient tools for consumers to complete a variety of financial transactions. Consumers use these electronic tokens to pay for daily expenses such as food, lodging, travel, and gasoline. The travel and tourism industry benefit greatly from the use of credit cards because in order to make hotel or rental car reservations, a consumer must provide a valid credit card (or in some instances a debit card) to reserve the desired hotel room or motor vehicle.\nThe credit card industry has grown considerably as electronic commerce and technology continues to advance. In the United States alone, there are millions of credit cards in circulation and the use of these cards will continue to grow as technology advances. With the advent of the internet, consumers have the ability to make purchases 24 hours a day from anywhere in the world by purchasing items through a retailer's website or toll free number.\nThe use of credit cards is particularly prevalent in the United States. Among the users of credit cards are those who are visually impaired. However, the use of credit cards sometimes requires that a consumer read the credit card account numbers or verification numbers from the card. Reading these account numbers can be inconvenient or even impossible for visually impaired individuals. Over ten million visually impaired card holders cannot read account numbers or the verification number for transaction validation purposes.\nPresently known techniques for permitting visually impaired individuals to read information, such as on a telephone or in an elevator, involves the use of Braille characters. Automated teller machines (ATM) and other terminals, public transit, elevators offer Braille pads to allow reading of instructions and other relevant information. In this way, the visually impaired individual may interpret the Braille characters on the keypad to facilitate dialing a phone number or accessing a particular floor in a hotel or an apartment. There is no known mechanism that permits the visually impaired individuals to read account and verification numbers on a credit card.\nBraille printers and other devices for generating Braille characters on printed documents are generally well known. Automated embossing systems have existed for a number of years. Primarily, these are devices for embossing letters or dots into a sheet of paper or plastic material, such as for embossing lettering into credit card blanks. Among these products are a number of Braille printers which use embossing heads which produce patterns of embossed dimples or dots corresponding to the Braille alphabet. Embossing creates raised elements, or bumps, that project from one surface of an object, such as a card, and a mirror image pattern of indentations on an opposite surface of the object. Embossing a set of characters in a region on one surface precludes embossing another set of characters on the contralateral region, i.e. the corresponding region on the opposite surface, as embossing the later formed set of characters typically obliterates part of the previously formed set.\nCredit card issuers must comply with standards promulgated by the International Organization for Standardization (ISO) in Geneva, Switzerland with respect to the dimensions of the credit card and the placement of characters, indicia, and alternative technology on cards. Such prior embossing and imprinting devices have not, however, been suitable for producing Braille characters on credit and debit cards which can be read by a visually impaired person or that complies with the requirements of ISO 7811-3 and ISO 7813 standards for financial transaction cards.\nVisually impaired card holders are limited to after-market solutions, such as Braille tape and embossers. These after market solutions require that an adhesive be placed on a tape or other printed material. However, these adhesive tactile solutions violate the ISO 7811-3 and ISO 7813 standards, requiring that alternative technologies can not be placed in areas of the card where they could adversely affect card features, such as the magnetic stripe. Attaching any type of adhesive on a credit card risks depositing adhesive materials in the card readers. In addition, the adhered indicia is easily removed through normal wear and tear by the consumer."}
{"text": "Typically, in a computer network, when a client requests content from a server, the request is sent to the server through a number of intermediary devices, each of which may alter the request in some way, according to rules installed on the intermediary device. The intent of these rules, and their embodiment, are commonly known as policy.\nManagement of the typical request/response process, however, is often complicated by a number of factors. First, there are a number of hardware and software computing elements involved, each of which is affected by policy, thus adding complexity to the management of the process.\nSecond, information required for policy decisions arrives at different times. For example connection information, request information and response information form three discrete “bundles” of information that become available at different times and within different processing subsystems. Thus the typical process is difficult to manage from a timing standpoint.\nFinally, the typical process is often extended in time, during which the policy rules may change, thus increasing the possibility of conflicting policy versions.\nThese factors, and others, add unpredictability to the typical process. Therefore, arbitrary changes to in-progress traffic, either by the client, by the intermediary device, or by other network devices, may possibly result in unexpected results to the client or to the network. Unexpected results may lead to relatively benign problems, such as causing irritation to the network user, or may lead to quite serious problems, such as causing security breaches in the network.\nTherefore, there is a need for a mechanism to allow for the uniform application of policy across separate processing elements within an intermediary device."}
{"text": "1. Field of the Invention\nThis invention relates to the field of artificial dentures, and more particularly, to such dentures prepared from polyurethane elastomers.\n2. Description of the Prior Art\nIt has been proposed to provide dentures with a soft layer in contact with the gums and other mouth parts to provide tissue relief. Such soft layers have been composed of acrylics, silicones, and like rubber-like materials. But on aging, such soft layers tend to harden and give off undesirable odors. In addition, some decomposition of the polymer may also occur presumably due to an oxidation process as well as to pH fluctuations within the mouth. By way of overcoming these disadvantages, U.S. Pat. Nos. 4,024,636 and 4,080,412, both to Colpitts et al., and both incorporated by reference herein, describe dentures in which teeth are anchored in a gum member comprising a tooth-holding portion fabricated from a hard nonhydrophilic polyurethane elastomer having a hardness of not less than about Shore D40, and a mouth-engaging portion fabricated from a soft non-hydrophilic polyurethane elastomer having a hardness of not greater than Shore A65 integrally and chemically bonded into a unitary mass. U.S. Pat. No. 4,024,637 to Colpitts which is also incorporated by reference herein describes a denture in which hard non-hydrophilic polyurethane elastomer teeth are imbedded in and chemically bonded to a soft non-hydrophilic polyurethane elastomer. Preferred non-hydrophilic elastomers are those formed by isocyanate-terminated prepolymers which are cross-linked or cured by mixing with a cross-linking agent and heating as required to effect curing. Isocyanate-terminated prepolymers suitable for preparing the hard non-hydrophilic polyurethane elastomers are prepared by the reaction of polyether diols or triols with aliphatic or cycloaliphatic or aralkyl di- or polyisocyanates in proportion to give free NCO groups. The prepolymers are then cured or cross-linked with a diol, polyol, an alkanolamine, a diamine or a tertiary amine containing polyol, or blends thereof. Advantageously, the diol or polyol is a polyether diol or polyol or a hydroxyl-terminated prepolymer. While dentures prepared from the foregoing polyurethane resins represent a significant advance over dentures fabricated from such materials as acrylic resins, silicones, and the like, it has been observed that the polyurethanes are susceptible to distortion under the conditions prevailing in the mouth. Such distortion, accompanied by a loss in dimensional stability of the dentures, interferes with the ability of the dentures to retain a good fit with the mouth and remain in place without slippage."}
{"text": "Flexible pipe is useful in a myriad of environments, including in the oil and gas industry. Flexible pipe may be durable and operational in harsh operating conditions and can accommodate high pressures and temperatures. Flexible pipe may be bundled and arranged into one or more coils to facilitate transporting and using the pipe.\nCoils of pipe may be positioned in an “eye to the side” or “eye to the sky” orientation. When the flexible pipe is coiled and is disposed with its interior channel facing upwards, such that the coil is in a horizontal orientation, then the coils of pipe are referred to as being in an “eye to the sky” orientation. If, instead, the flexible pipe is coiled and disposed such that the interior channel is not facing upwards, such that the coil is in an upright or vertical orientation, then the coils of pipe are referred to as being in an “eye to the side” orientation.\nThe flexible pipe may be transported as coils to various sites for deployment (also referred to as uncoiling or unspooling). Different types of devices and vehicles are currently used for loading and transporting coils of pipe, but usually extra equipment and human manual labor is also involved in the process of loading or unloading such coils for transportation and/or deployment. Such coils of pipe are often quite large and heavy. Accordingly, there exists a need for an improved method and apparatus for loading and unloading coils of pipe."}
{"text": "Many men and women suffer from urinary incontinence, to varying degrees, and for various reasons. Whether their particular case is extreme or mild, of long duration or of recent onset, constant or discontinuous, and regardless of the cause, urinary incontinence may pose a significant embarrassing and/or uncomfortable health related challenge. As such, mitigation of urinary incontinence, and of the discomfort and visual impact of urinary incontinence, has long been the goal of many different apparatus and medically related techniques. Indeed, there are numerous different types of incontinence pads (including but not limited to pads, guards, pouches, briefs, and diapers) available for men and women. However, while some conventional pads may indeed effectively capture and retain leaked urine, conventional pads, in focusing only on capture of urinary emissions, do not properly address mitigation of the incontinence itself, i.e., reducing the frequency of emissions and/or reducing the amount of urine leaked during an incontinence event. Even as to those pads that do effectively capture and retain leaked urine, there are still improvements that can be made. Particular shortcomings of conventional pads relative to fluid retention may relate, for example, to fluid spill-out (i.e., leakage of fluid out of the pad along the edges of the pad) and uncomfortable, perhaps noisingly embarrassing fluid slosh (from leaf to right, or vice versa, within the pad), which often may occur while the user is sitting, or moving during sitting. And, perhaps more significantly, as mentioned, conventional pads tend not to provide any mechanism for mitigating the incontinence itself. It is of note that the incontinence mitigation that is achieved via certain embodiments of the inventive technology disclosed herein may have conventionally been achieved only via surgery. As such, particular embodiments of the inventive technology may achieve a desired degree of incontinence mitigation via measures—wearing of an inventive pad—that are much less drastic and interfering than surgery."}
{"text": "The present invention relates to an image forming apparatus such as a copying machine.\nRecent image forming apparatuses such as copying machines have various functions. Parameters which represent such functions as the copying number, the copying paper size, the copying magnification, the copying density and the like will be referred to as image forming instructions hereafter. When such a copying machine is manufactured or is placed and used in an office, it may be necessary to check all the functions of the copying machine. A conventional checking operation is performed in such a manner that after all the functions of the copying machine are manually set at all possible values, a copying start instruction (depression of a start key) is provided. For this reason, the checking operation becomes complex and time-consuming according to the various functions of the copying machine."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor devices and image forming apparatus, and more particularly to a semiconductor device having a light-emitting element array adaptable for use as a light source in electronic equipment including but not limited to laser printers and laser display units plus optical communications systems as well as optical signal processing apparatus, along with an image formation apparatus using the semiconductor device.\n2. Description of the Related Art\nIn image formation equipment such as laser printers or the like, in the case where a light-emitting element array with a plurality of light-emitting elements is employed as the light source for producing light according to respective pixels of an image to be formed, it is required that each light-emitting element be driven in a way independent of others; to this end, one typical prior known approach is to provide lead wires for electrically feeding respective light-emitting elements (referred to as \"individual leads\" hereinafter) on a one-per-element basis. Accordingly, in the case of a light-emitting element array constituted from n rows and m columns of light-emitting elements, i.e. n.times.m light-emitting elements, a great number of connection leads as many as n.times.m lines are required resulting in creation of a problem as to an increase in production cost while making it difficult to achieve higher integration density of light-emitting elements.\nOne approach to avoiding this problem has been described in Published Unexamined Japanese Patent Application (\"PUJPA\") No. 61-31271, which teaches the use of a technique for employing an LED array light source to form a hard copy or \"print\" image of high density. As shown in FIG. 11, the LED array 29 is made up of LED elements laid out two-dimensionally which are associated with a mesh-like or \"matrix\" pattern of rows of parallel leads (referred to as \"cathode leads\" hereinafter) 3 and columns of leads (\"anode leads\") 4. These cathode leads 3 and anode leads 4 have certain terminate ends at which wire bonding pads 33 and 32 are provided for use in individually driving the LED elements.\nAnother related art technique is found in PUJPA No. 7-503104 (\"JP-A-7-503104\"), which suggests use of a lamination of vertical-cavity surface-emitting laser structures with a transistor structure to thereby form a matrix lead pattern associated with respective transistor elements in the transistor structure. Driving a given transistor element permits light production at a corresponding vertical-cavity surface-emitting laser element that is placed on the upper side or the lower side of the transistor being presently driven.\nThe related art techniques set forth in the Japanese documents JP-A-61-31271 and JP-A-7-503104 require the use of a reduced number (n+m) in total of leads for the light-emitting element array of a 2D matrix of n rows and m columns of light-emitting elements; thus, it becomes possible to greatly reduce the required number of connection leads when compared to the case of the individual lead connection stated above. As a single wire bonding is formed per lead for example, a decrease in number of connection leads would lead to a decrease in number of wire bonding and also to a decrease in channel number of driver circuitry for use in driving the light-emitting element array. This in turn makes it possible to reduce production costs of the entire of equipment that employs this light-emitting element array, and also to achieve high integration density of such light-emitting element array.\nUnfortunately, the related art approaches as taught by the Japanese documents JP-A-61-31271 and JP-A-7-503104 cited above are encountered with a problem which follows. In case the light-emitting elements constituting the light-emitting element array are further increased in number and disposed with higher integration density while narrowing the interval between rows of light-emitting elements and/or between columns thereof, the layout pitch of such leads (anode leads and cathode leads) becomes narrower resulting in incapability of implementation of the intended wire bonding, which in turn makes it impossible to achieve electrical interconnection required. In other words, the minimum lead pitch enabling wire bonding or else limits the row distance and column distance of those light-emitting elements making up a light-emitting element array. This will be explained more practically below.\nIn currently available wire-bonding techniques the minimum allowable pitch is 80 .mu.m, or more or less. Suppose that wire-bonding pads are arrayed in a staggered or \"zigzag\" pattern in which pads are arrayed in two parallel lines with each pad in one line opposing a midway between adjacent one of the pads forming the other line. If this is the case, even where the minimum pitch of wire-bonding pads in each line is set at approximately 80 .mu.m, the minimum lead pitch can be reduced down to about 40 .mu.m, which is half of the minimum pad pitch. However, even in this case, the minimum lead pitch stays as large as about 40 .mu.m.\nAlternatively, imagine that the required number of connection leads is relatively less. In this case, expanding those portions of leads near the wire bonding pads into a sector shape might enable enlargement of the wire-bonding pad pitch. However, if the required number of such leads increases beyond several hundreds, then this pad pitch enlargement approach will no longer work well. This can be said because an increase in pitch of wire-bonding pads would result in an excessive increase in size of a substrate used. In view of this, the minimum layout pitch of those light-emitting elements constituting a light-emitting element array of relatively high density has been typically limited to 40 mm, or more or less.\nIt has been discussed that the related art approaches stated above suffer from a problem that a light-emitting element array having multiple light-emitting elements is strictly required to set the pitch of its leads at a limited value that permits electrical connection by wire bonding or the like, which hardly allows the pitch to decrease below a prespecified value of several tens of .mu.m. This is a serious bar to reduction of the row distance and column distance of respective light-emitting elements, which in turn makes it impossible, or at least greatly difficult, to manufacture the intended light-emitting element array with high integration density required."}
{"text": "Slurry phase polymerisation of olefins is well known wherein an olefin monomer and optionally olefin comonomer are polymerised in the presence of a catalyst in a diluent in which the solid polymer product is suspended and transported.\nThis invention is specifically related to polymerisation in a loop reactor where the slurry is circulated in the reactor typically by means of a pump or agitator. Liquid full loop reactors are particularly well known in the art and are described for example in U.S. Pat. Nos. 3,152,872, 3,242,150 and 4,613,484.\nPolymerisation is typically carried out at temperatures in the range 50-125° C. and at pressures in the range 1-100 bara. The catalyst used can be any catalyst typically used for olefin polymerisation such as chromium oxide, Ziegler-Natta or metallocene-type catalysts. The product slurry comprising polymer and diluent, and in most cases catalyst, olefin monomer and comonomer can be discharged intermittently or continuously, optionally using concentrating devices such as hydrocyclones or settling legs to minimise the quantity of fluids withdrawn with the polymer.\nThe loop reactor is of a continuous tubular construction comprising at least two, for example four, vertical sections and at least two, for example four, horizontal sections. The heat of polymerisation is typically removed using indirect exchange with a cooling medium, preferably water, in jackets surrounding at least part of the tubular loop reactor. The volume of the loop reactor can vary but is typically in the range 20 to 120 m3; the loop reactors of the present invention are of this generic type.\nMaximum commercial scale plant capacities have increased steadily over the years. Growing operating experience over the last few decades has led to operation of increasingly high slurry and monomer concentrations in reaction loops. The increase in slurry concentrations has typically been achieved with increased circulation velocities achieved for example by higher reactor circulation pump head or multiple circulation pumps as illustrated by EP 432555 and EP 891990. The increase in solids loading is desirable to increase reactor residence time for a fixed reactor volume and also to reduce downstream diluent treatment and recycling requirements. The increased velocity and head requirement of the loop has however led to increasing pump design sizes and complexity, and energy consumptions as slurry concentrations increase. This has both capital and operating cost implications.\nHistorically the circulation velocity in the reaction loop has typically been maximised to ensure maintenance of good thermal, compositional and particle distribution across the reactor cross-section, particularly the avoidance of solids settling, stable flow characteristics, or excessive solids concentrations at the pipe wall rather than reduced to minimise pressure drop/power in the polymerisation loop.\nInadequate cross-sectional distribution could lead to increased fouling, reduced heat transfer and reduced polymer productivity and homogeneity. Construction and commissioning of new commercial plants is very expensive and therefore new designs seek to avoid or minimise changes to operating parameters that are seen to increase risk to the successful operation of the new unit."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a field emission device (FED), and more particularly, to a method of manufacturing a FED using a mask for performing multiple exposure to light and multi-patterning processes.\n2. Description of the Related Art\nIn general, thin films or thick films forming stacked structures for electronic devices are processed using physical, chemical, or physical-chemical methods. Here, masks for patterning processing target materials to have desired shapes are used.\nIn general, a mask is directly deposited to a processing target material, and is formed with a photoresist of a light-hardened or light-softened polymer. According to a conventional method, a conventional mask is formed through a series of processes such as photoresist coating, soft baking, exposing, developing, and hard baking, and the mask cannot be altered when these processes are finished. Accordingly, the conventional mask is used only once, and thus a new mask needs to be formed if the stacked structure requires an additional etching process for forming a new pattern.\nSince the conventional photolithography process includes many operational units, they should be simplified when considering manufacturing costs and productivity."}
{"text": "1. Field of the Invention\nThe present invention relates to a self-oscillating full-bridge driver IC, and more particularly to a driver having a pre-charging circuit, which may include a bootstrap circuit for pre-charging a bootstrap capacitor before starting up or transferring signals from the oscillator.\n2. Related Art\nA self-oscillating full-bridge driver of background interest is the IR2153 integrated circuit. A functional block diagram of the IR2153 is shown in FIG. 1. The IR2153 has several features of particular interest:\nUnder voltage lockout (UVLO) is the name for a protection feature built into several control IC's manufactured by International Rectifier. This block, marked as UV Detect on FIG. 1, ensures that the gate drive outputs, HO and LO are both low should bias become too marginal for comfortable gate drive to the output transistors. The UVLO circuit also assures a repeatable start-up sequence and controls bias current needed for various elements of the IC. The UVLO circuit monitors the bias supply Vcc.\nThe UVLO circuit is further provided with hysteresis to improve noise immunity. Instead of switching on/off at one voltage level, two threshold levels are defined; one sensitive to a rising edge and the other to a falling edge. In product data sheets describing Vcc undervoltage lockout, these two thresholds are listed as VccUV+ and VccUV−. Hysteresis is a measure of the difference between these two thresholds and is separately listed on data sheets as VccUVH. VccUV+ is the threshold at which the chip is first enabled, when Vcc is rising. Once enabled, Vcc is allowed to drop slightly without adverse effects.\nThe IR2153 also has micropower start-up, which means that the bias current needed when the IC is in a quiescent state (i.e. not oscillating and Vcc<VccUV) is reduced to a very low level. Micropower start-up is beneficial because off-line applications usually have a single resistor from a high voltage bus for start-up bias to keep component cost and complexity low. In some cases this resistor provides bias for all modes of operation(both start-up and running, where current demand is highest). However, an auxiliary low voltage source is more common for this purpose.\n“Bootstrapping” is beneficial because an auxiliary supply can be derived with much higher efficiency than from a simple dropping resistor connected to the line. In bootstrap start-up schemes, the micropower feature allows the pull-up resistor to be increased in value compared to schemes in which a single dropping resistor provides all the bias. This is because the resistor needs only to overcome the quiescent bias current requirements for the IC and raise Vcc to UV+. When the IC starts oscillating, the additional bias requirement is met by an auxiliary supply.\nThe term “bootstrap” is also used to describe the diode and Vbs capacitor arrangement used in the typical circuit configuration. This part of the circuit serves a different function than the resistor in a bootstrap “start up” scheme. Bootstrap diode/capacitor circuits are described in Design Tip DT-2, by Jonathan Adams, which is available at www.irf.com. incorporated by reference.\nThe IR2153 further incorporates a fast shutdown mode providing a protective action against fault conditions that would otherwise destroy output switches. Fusing is the simplest option but is often ineffective at preventing damage to semiconductors and is not self-resetting, unless comparatively more expensive devices are used. If the load is capacitively coupled or resonant in nature, for example as in most electronic ballast designs, ample protection may be afforded by merely halting the oscillator and preventing further switching cycles. However, adequate protection sometimes calls for both FET's to be rapidly turned off.\nOne way to implement shutdown is to pull Vcc below VccUV− so as to turn off both FETs. This requirement may to an extent conflict with the need to provide a healthy charge reservoir and solid decoupling. The shutdown network must be capable of rapidly pulling charge from the Vcc reservoir capacitor (and the decoupling capacitor, if present). Low holding current thyristors can be used to crowbar Vcc and latch the chip off until the supply is recycled. This approach is both simple and inexpensive but does not support automatic power-on reset.\nThe IR2153 offers a convenient shutdown solution by adding a second function to the Ct pin. A third functional threshold is added to Ct and set at Vcc/6. Below this threshold a fast shutdown mode is invoked and both output buffers are set low with minimal delay. This allows a simple open-collector NPN transistor or similar to be used as part of an inexpensive protection scheme.\nIR2153 can be if desired be driven directly as a slave device, for example to complete a fill bridge. The Ct input can be directly fed from a square wave source, provided Vct high>⅔ Vcc and Vct low<⅓ Vcc. The Rt output of any self-oscillating control IC can generate a master clock to feed directly into the Ct input of another, effectively bypassing the internal oscillator of the slaved device. This can be a useful way to implement a full bridge using just one type of IC or hybrid.\nIR2153 further has improved dead-time accuracy with zero average temperature coefficient. Dead time is the period during which both HO and LO outputs are intentionally low. This period is fixed inside the IC and serves several essential functions. The primary function of dead-time is to prevent cross conduction or shoot-through in half-bridge designs. If the load is resonant in nature, such as in electronic ballasts and resonant mode power supplies, dead-time also helps maintain zero voltage switching (ZVS). Sometimes called soft switching, this technique significantly reduces switching losses.\nCross-conduction will occur in half bridge circuits if both high and low side transistors are either fully or partially on at the same time. The resulting short across the supply exacerbates EMI, increases dissipation and may destroy power switches, the control IC, or both if extreme. MOSFET turn on/off times are often inequal and vary in production, so dead-time offers a guard band to account for these differences. To prevent cross conduction, or shoot-through, the IR2153 is available with various dead time options. Moreover, switching times can be easily modified using small signal diodes across series resistors in the gate drive loop if necessary.\nWhen the load is inductive or resonant, circuit efficiency is generally highest when zero voltage switching (ZVS or soft switching) occurs. In ZVS, there is just enough time when both FETs are off (approximately equal to the dead-time) to allow load energy to swing the output to the opposite rail, where current will flow into the FET body drain diode. This desirable process is called self-commutation.\nSoft switching (ZVS) is more efficient than hard switching. MOSFETSs become hot when the load is removed.\nSelf-commuation is almost completely lossless because output capacitance of the half bridge rings with the inductive component of the load when both half bridge transistors are off. When the next transistor turns on, the voltage across it is already at or close to zero, so its internal capacitance is already discharged. When its load is not resonant or inductive in nature, the output cannot self-commute to the opposite rail during off time, so the voltage across the next transistor to turn on is high. It must therefore discharge its own output capacitance in addition to providing load current. This explains why half bridge transistors often get hot when there is no load at all, or in ballasts when the lamp is removed and the oscillator is allowed to continue running.\nIf ZVS is required but the dead-time is too short, there is insufficient time for the half-bridge output voltage to completely self-commutate to the opposite rail. This leaves a fraction of the bus voltage remaining across the output capacitance of both transistors which the next transistor to switch on must discharge. This is partial hard switching (or partial soft switching). Conversely, if the dead-time is too long, the output voltage will self-commutate to the opposite rail, but a short time later, the load current will reverse or ring, producing a voltage transient. Again, the results is partial hard switching and less than optimal switching efficiency. For ZVS, optimal dead time occurs when there is just enough time for self commutation to occur. IR2153 is available with several dead-time options to help meet this requirement. If further adjustments are necessary, small signal diodes can be used in the gate drive loop to independently change charge/discharge timing.\nIR2153 output buffers also have a reduced di/dt output stage. This means that the peak rate of change of current (rate of current rise and fall for turn-on and turn-off) has been deliberately reduced. In effect the IR2153 output buffers are turned on and off at a controlled rate.\nBenefits of reduced di/dt output buffers include lower gate ringing (Vgs overshoot/undershoot) from a combination of di/dt and gate circuit parasitics. In some cases gate resistors can be eliminated altogether, however this depends on several factors such as layout and transistor type used. Rf interference may also be reduced."}
{"text": "The use of chat or instant messaging communications in a networked computer environment is well known. The America Online™ Instant Messaging system (AIM) is one well known system for controlling instant messaging in a networked environment. In these prior art systems, two computer users can communicate non-persistent text messages in real time using an instant message (IM) client on their computers in concert with an IM server.\nMost messaging services are subscription-based or user-identity-based and may generate large numbers of content followers or users of particular message or content sources (denoted herein as subscribers). These content followers or subscribers can form communities or social networks around a particular content source or content distribution system. Social networks have gained in popularity as people have used messaging as a basis for connecting with each other.\nAs the numbers and size of the user pool, subscribers, and social networks expand, it becomes more difficult to track and manage the subscribers, the listening users, and the degree to which the users are involved with the message content. Similarly, it becomes more difficult to identify and rank the most popular content items being consumed across a variety of content sources and social networks.\nThus, a system and method managing multiple queues of non-persistent messages in a networked environment are needed."}
{"text": "This invention relates to liquid level sensors, typically the sending unit for a fuel gauge. More particularly, this invention relates to liquid level sensors that employ an optical waveguide that is immersed in the fluid to provide a continuous measurement of the liquid level.\nThe need for continuous liquid-level measurements exists in numerous commercial and military applications, such as in the fuel, oil and water tanks of aircraft, automobiles and trucks. Less mobile applications include storage tanks or fuel dispensing, waste water, home, chemical, and food processing purposes, to name but a few. Electrical sensors are of particular concern in many cases, particularly those involving flammable fluids. Hazards are apparent from electric sparks from such sensors in these potentially explosive environments. In other environments, electromagnetic interference may overwhelm the signals generated by these sensors. Conventionally, the level of a liquid in a vessel is detected using a float with a mechanical or magnetic coupling to an external gauge, an ultrasonic or optical transducer which measures time of flight to deduce the liquid level, or a parallel wire capacitance sensor that monitors the change in the dielectric constant between the wires associated with a change in liquid level. Capacitance sensors lose accuracy with increasing amounts of water in the fuel. Sensing of fluid level by visual inspection is only viable in limited circumstances since most usage's require remote and continuous sensing of fluid levels.\nA number of fluid level sensors based on fiber optics have been proposed. They are attractive because they introduce no electrical energy into the tank, are insensitive to electromagnetic interference, have no moving parts and can provide a continuous measurement of fluid level. A number of older fiber optic fluid level sensors are described in the article \"Experimental Investigations on Fibre Optic Liquid Level Sensors and Refractometers,\" by K. Spenner et al. in IEE OFS 221 pages 96-99. These implementations employ two separate optical fibers and are primarily for the measurement of discrete levels in the vessel, e.g. to sense whether the vessel is completely full or completely empty or at some intermediate point. It would be more desirable to have a fiber optic sensor that could provide continuous measurements.\nAn early continuous measurement fiber optic sensor is described in \"Fibre Optic Fluid Level Sensor,\" by M. Belderdid, N. Ghanderharioun, and B. Brennan in the Proceedings of the SPIE Conference 566 Fibre Optic and Laser Sensors IlIl (1985) pages 153-158. This sensor is based on the bending or cladding loss principle, consisting of large loops of a single fibre, the loops being of ever increasing diameter and suspended in the liquid. U.S. Pat. No. 4,870,292 to Alpert et al., teaches a fluorescent doped detector fiber to collect light reflected from a source fiber in the presence of air; the light is refracted away when a fluid of higher refractive index is present. The output signal from the fluorescent doped fiber is analogous to the fluid level.\nU.S. Pat. No. 4,942,306 builds upon the '292 patent above by embedding one end of the source optical fiber into a transparent substrate such that the injected light enters the substrate at such an angle that it will refract out of the substrate when the substrate is in the liquid but will continue to be totally internally reflected when the substrate is in the air above the liquid. The detecting fiber is again a separate fluorescent doped fiber that is placed against the transparent substrate, typically being wrapped in grooves around the outside of the substrate.\nNevertheless, there remains a need in the art for a simple, low cost liquid level sensor based on this fiber optic technology, preferably realized as a single optical element immersed in the fluid."}
{"text": "The field of the invention is static control containers and the invention relates more particularly to containers for transporting printed circuit boards and other devices which could be harmed by electrostatic discharge, electromagnetic interference or radio frequency interference as well as physical damage resulting from shock and vibration.\nIn the fabrication of modern electronic equipment, printed circuit boards are often fabricated in one location for use at a remote site. Although such boards can be fabricated with great concern for a static-free environment, the shipping step has, in the past, led to damage or destruction of the board as a result of electrostatic discharge. Such boards are also often physically fragile and the boards should be transported in such a way where they will not bend or be subject to impact. Plastic bags fabricated from anti-static material and including a conducting layer such as aluminum wire mesh screen have been used for transporting, but such bags do not provide protection against physical impact and, also, are typically relatively fragile and are usually not reused. Because of this one-time use, the cost of such packs is relatively high. Corrugated containers, including a conductive metal foil layer, have been used, but such containers are also typically not reusable, although they do provide improved protection from physical damage as compared to a pouch."}
{"text": "Small electronic devices are often portable in nature and, as such, require a portable source of power such as a battery. A battery, which typically includes a battery housing, may be conveniently attached to an electronic device where electrical contact is made through conventional electronic coupling.\nSeveral mechanisms have been devised to enable a user to easily connect and disconnect a battery from an electronic device. One such mechanism has been a latching mechanism. As an electronic device becomes smaller, however, volume within the electronic device cannot afford a latching mechanism which is extensive and complex. Thus, it would be advantageous for a latching mechanism to be placed in locations other than within the electronic device itself.\nIn addition, an electronic device having a battery may be subject to some mechanical stresses, such as those stresses which occur when an electronic device is dropped accidentally. Depending on the latching mechanism's design, the latching mechanism may break or wear from such stresses. Wearing the latching mechanism may cause the battery to be loosely attached to the electronic device. At best, a loosely attached battery may rattle on the electronic device; at worst, the battery may lose electrical contact with the electronic device. Therefore, it would be advantageous to provide a latching mechanism that is durable. In addition, a battery having a simple design for manufacturing is often desirable for reducing manufacturing effort, cost, and time. Also, a battery having a light weight is desirable when considering product cost and electronic device portability."}
{"text": "This invention relates to a medical device that is useful in aiding physicians to make complex multi-plane incisions in a patient. More particularly this invention relates to a form that distorts a portion of a patient's body in such a way that allows the physician to make a straight line planar incision in the distorted portion of the patient's body to form a complex multi-plane incision in the patient's body when the form is removed from the patient's body.\nIn various surgical procedures, the physician typically has to make an incision in the patient in order to remove unwanted tissue, repair damaged tissue, or implant a device to improve the patient's well being. In certain cases, all three of these activities or a combination thereof must be done in a single procedure. For example, in cataract surgery, the physician removes the natural lens that has been clouded by a cataract from the patient's eye and replaces it with an artificial lens that will improve the patient's eyesight. In order to perform this procedure, an incision is made in the cornea or sclera of the eye. This provides the physician with access to the patient's lens. The clouded lens is cut loose and removed. There are a number of different procedures that are used to remove a patient's lens that has a cataract. Two of the more common techniques are known as extracapsular surgery and phacoemulsification. In extracapsular surgery, the physician removes the lens leaving behind the posterior capsule. In phacoemulsification, the physician fragments the lens by ultrasonic vibrations. The lens is simultaneously irrigated and aspirated. After the lens is removed, the physician then inserts an artificial lens known as an intraocular lens (IOL) into the eye behind the iris. Two tiny C-shaped arms connected to the IOL eventually become scarred into the side of the eye and hold the IOL firmly in place.\nAlthough cataract surgery is considered routine, it is not foolproof. For example, the incision through which the devices that are used to perform the procedure are inserted must provide an opening that is substantially the same circumference as the probes and other surgical devices that are used to perform the procedure. This minimizes trauma to the eye and facilitates rapid healing of the eye after the procedure. In addition, in the case of phacoemulsification, if the incision is too small, corneal tissue surrounding the incision may contact the ultrasound probe and be burned. Alternatively, if the incision is too large, leakage from the eye after the procedure may occur causing iris prolapse and subsequently endothelial cell loss.\nAnother common problem with cataract surgery is suture induced astigmatism. Since an incision is made in the eye to perform the procedure, some mechanism must be used to ensure that the incision remains closed during the healing process to prevent a path for infection into the eye and so that the eye heals properly. In the past, sutures have been used to close the incision. However, the use of sutures in the eye to seal the incision after the cataract surgery may alter the shape of the eye. Such an alteration of the shape of the eye may result in astigmatism. In addition, the use of sutures may cause other complications such as eye irritation, suture abscesses, suture extrusion and foreign body reaction,\nIn an effort to avoid the use of sutures in cataract surgery, new sutureless techniques have been developed. These sutureless techniques involve the same standard procedures used to remove the cataract and implant the IOL but require the physician to make an incision in the eye having a particular shape or geometry. By making an incision having a particular geometry, the normal internal pressure in the eye pushes against the eye at the location of the incision to keep the incision closed.\nOne particular sutureless technique that is currently being used is known as the clear cornea approach and requires a particular incision that is self-sealing to be made in the eye. The shape of one particular incision comprises a pair of contiguous circular segments that have their centers of radius on opposite sides of the incision, i.e. the incision has a substantially S shaped geometry. Unfortunately, because of the complex, multi-plane geometry of the incision, it is difficult to form consistently. This problem is especially acute for new physicians and physicians who do not perform this technique often enough to maintain their proficiency. If such incisions are improperly formed, a number of complications can occur. For example, endophthalmitis, induced astigmatism, damage to the cornea, the iris or other intraocular tissue can occur, and impaired intraoperative vision of the physician can result."}
{"text": "The present invention relates in general to semiconductor devices and to methods of fabricating semiconductor devices, and more particularly, to methods of fabricating cross-coupled line segments for use, for instance, in fabricating of one or more semiconductor devices.\nAs is known, semiconductor device fabrication typically involves the process of transferring a mask pattern to a wafer using an etching to remove unwanted material, for instance, to facilitate forming fin structures and/or gate structures of fin-type field-effect transistors (FinFETs) or to facilitate forming field-effect transistors (FETs) in general. This process is generally referred to as lithographic processing. As the size of technology nodes continues to decrease, significant challenges continue to arise due (in part) to issues related to traditional lithographic processing techniques, including issues related to mask formation."}
{"text": "At initialization of conventional multi-board systems, each circuit board attempts to locate its memory, interrupt addresses, registers, bus grants, etc., at the same address. A number of jumpers or strapping devices are used to configure or memory map the memory, interrupt addresses and registers so that they do not overlap. However, most conventional systems require a qualified individual who can reconfigure the circuit boards to place the jumpers or strapping devices in the proper position and/or location so that each circuit board is configured in a particular order.\nWhen plugging a mezzanine (or daughter) circuit board onto either side of a master circuit board (or a mother board), enable signals are sent in a daisy chain fashion from one circuit board to another circuit board. Therefore, if a mezzanine board is placed on top of the master board, the enable signals are passed from the master board to the mezzanine board. If the mezzanine board is placed on the other side of the master board, the enable signals are passed in that direction from the master board to the mezzanine board. However, a skilled person is needed who knows how to manually change the jumpers or strapping devices on both the master and mezzanine boards so that the enable signals travel in the proper direction.\nIf mezzanine boards are to be placed on both sides of the master board, the jumpers or strapping devices will have to be changed by a skilled person. In other words, each time the master and mezzanine boards are rearranged, jumpers or strapping devices are required so that the entire system initializes and configures correctly and functions properly. Thus, there is a significant need for a multi-board system which eliminates the intervention of a skilled person for changing jumpers or strapping devices whenever the master and mezzanine boards are rearranged."}
{"text": "The present invention provides compounds, pharmaceutical compositions containing one or more of those compounds or their pharmaceutically acceptable salts, which are effective in inhibiting the binding of various chemokines, such as MIP-1α, leukotactin, MPIF-1 and RANTES, to the CCR1 receptor. As antagonists or modulators for the CCR1 receptor, the compounds and compositions have utility in treating inflammatory and immune disorder conditions and diseases.\nHuman health depends on the body's ability to detect and destroy foreign pathogens that might otherwise take valuable resources from the individual and/or induce illness. The immune system, which comprises leukocytes (white blood cells (WBCs): T and B lymphocytes, monocytes, macrophages granulocytes, NK cell, mast cells, dendritic cell, and immune derived cells (for example, osteoclasts)), lymphoid tissues and lymphoid vessels, is the body's defense system. To combat infection, white blood cells circulate throughout the body to detect pathogens. Once a pathogen is detected, innate immune cells and cytotoxic T cells in particular are recruited to the infection site to destroy the pathogen. Chemokines act as molecular beacons for the recruitment and activation of immune cells, such as lymphocytes, monocytes and granulocytes, identifying sites where pathogens exist.\nDespite the immune system's regulation of pathogens, certain inappropriate chemokine signaling can develop and has been attributed to triggering or sustaining inflammatory disorders, such as rheumatoid arthritis, multiple sclerosis and others. For example, in rheumatoid arthritis, unregulated chemokine accumulation in bone joints attracts and activates infiltrating macrophages and T-cells. The activities of these cells induce synovial cell proliferation that leads, at least in part, to inflammation and eventual bone and cartilage loss (see, DeVries, M. E., et al., Semin Immunol 11(2):95-104 (1999)). A hallmark of some demyelinating diseases such as multiple sclerosis is the chemokine-mediated monocyte/macrophage and T cell recruitment to the central nervous system (see, Kennedy, et al., J. Clin. Immunol. 19(5):273-279 (1999)). Chemokine recruitment of destructive WBCs to transplants has been implicated in their subsequent rejection. See, DeVries, M. E., et al., ibid. Because chemokines play pivotal roles in inflammation and lymphocyte development, the ability to specifically manipulate their activity has enormous impact on ameliorating and halting diseases that currently have no satisfactory treatment. In addition, transplant rejection may be minimized without the generalized and complicating effects of costly immunosuppressive pharmaceuticals.\nChemokines, a group of greater than 40 small peptides (7-10 kD), ligate receptors expressed primarily on WBCs or immune derived cells, and signal through G-protein-coupled signaling cascades to mediate their chemoattractant and chemostimulant functions. Receptors may bind more than one ligand; for example, the receptor CCR1 ligates RANTES (regulated on activation normal Tcell expressed), MIP-1α (macrophage inflammatory protein), MPIF-1/CKβ8, and Leukotactin chemokines (among others with lesser affinities). To date, 24 chemokine receptors are known. The sheer number of chemokines, multiple ligand binding receptors, and different receptor profiles on immune cells allow for tightly controlled and specific immune responses. See, Rossi, et al., Ann. Rev. Immunol. 18(1):217-242 (2000). Chemokine activity can be controlled through the modulation of their corresponding receptors, treating related inflammatory and immunological diseases and enabling organ and tissue transplants.\nThe receptor CCR1 and its chemokine ligands, including, for example MIP-1α, MPIF-1/CKβ8, leukotactin and RANTES, represent significant therapeutic targets (see Saeki, et al., Current Pharmaceutical Design 9:1201-1208 (2003)) since they have been implicated in rheumatoid arthritis, transplant rejection (see, DeVries, M. E., et al., ibid.), and multiple sclerosis (see, Fischer, et al., J Neuroimmunol. 110(1-2):195-208 (2000); Izikson, et al., J. Exp. Med. 192(7):1075-1080 (2000); and Rottman, et al., Eur. J. Immunol. 30(8):2372-2377 (2000). In fact, function-blocking antibodies, modified chemokine receptor ligands and small organic compounds have been discovered, some of which have been successfully demonstrated to prevent or treat some chemokine-mediated diseases (reviewed in Rossi, et al., ibid.). Notably, in an experimental model of rheumatoid arthritis, disease development is diminished when a signaling-blocking, modified-RANTES ligand is administered (see Plater-Zyberk, et al., Immunol Lett. 57(1-3):117-120 (1997)). While function-blocking antibody and small peptide therapies are promising, they suffer from the perils of degradation, extremely short half-lives once administered, and prohibitive expense to develop and manufacture, characteristic of most proteins. Small organic compounds are preferable since they often have longer half lives in vivo, require fewer doses to be effective, can often be administered-orally, and are consequently less expensive. Some organic antagonists of CCR1 have been previously described (see, Hesselgesser, et al., J. Biol. Chem. 273(25): 15687-15692 (1998); Ng, et al., J. Med. Chem. 42(22):4680-4694 (1999); Liang, et al., J. Biol. Chem. 275(25):19000-19008 (2000); and Liang, et al., Eur. J. Pharmacol. 389(1):41-49 (2000)). In view of the effectiveness demonstrated for treatment of disease in animal models (see, Liang, et al., J. Biol. Chem. 275(25):19000-19008 (2000)), the search has continued to identify additional compounds that can be used in the treatment of diseases mediated by CCR1 signaling."}
{"text": "Persistent malware is a common type of malware that installs itself on a device and is thus able to perform actions whenever the device is on. So called bots (participating in botnets) can be classified as persistent malware that utilize a network and are widely held as one of the most serious threats on the Internet today. In general usage, the term “bot” refers to a computer infected with malware that allows the attacker to remotely control the computer to perform tasks on the behalf of the attacker, such as sending spam, stealing information, or launching attacks against other computers. A botnet is a collection of such bots under common control. Many techniques already exist for detecting malware infections on computers and other devices with similar capabilities, such as phones, tablets, etc. Existing host-based techniques include signature scanning for infected and manipulated files and behavioural monitoring for unusual and/or suspect processes on the device. Similarly, common network-based techniques rely on detecting signatures of attacks or malware communications or abnormal or suspect communication patterns, likely related to malware infections. Most techniques depend on capture and manual analysis of malware samples to create detection rules for that specific malware.\nExisting malware detection techniques to date known by the inventors are imperfect in terms of detection performance and their capabilities depend on where they are deployed, e.g., on the host or in the network. Since there is a constant arms-race between malware writers and the security industry trying to defend the users, there is constantly a need for new techniques to improve detection of threats such as persistent network-using malware."}
{"text": "The invention relates to a connector for making a connection underwater, sometimes known as a “wet-mate” or “wet-mateable” connector.\nSo called wet-mate connectors are used in underwater applications where it is necessary to make a connection, such as an electrical or optical connection, in an environment which is hostile to contact, for example in sea water, and which therefore requires special protection for the components that complete the connection.\nOne example of an application in which an electrical connection must be made in a harsh underwater environment is that of a well-head in a sub-sea oil well.\nAfter assembly of the well-head on the sea bed it is necessary to connect control cables to sensors and other electrical equipment associated with the well-head. The engagement of a control cable with a corresponding connector on the well-head may be carried out by a diver. The two connectable parts typically comprise a receptacle part and a plug part; the latter which becomes inserted within the former. Each part comprises a substantially cylindrical body part having within it the electrical contact. The electrical contact is typically provided with a protective apparatus to shield it from the surrounding sea water, in order to preserve the integrity of the connector and therefore the electrical connection when subsequently made.\nThe receptacle part houses a male connecting pin, and the plug part houses the complementary female contact socket. Each of the receptacle and plug is attached by a suitable termination means to respective electrical cables. In use, the receptacle part receives the plug part and as it does so the male pin contact penetrates and makes electrical connection with the female contact socket. Various designs of such connector exist in which there may be a single male pin engaging with a single contact module, or else a plurality of male pins and respective contact modules.\nIn all cases measures must be taken to prevent the electrical contacts from being exposed to sea water and other harmful matter, such as oil and drilling fluid for example. Maintaining a good seal around the electrical contacts may be necessary for long periods. In order to provide protection for the electrical contacts a number of mechanisms are employed. These include one or more wiper seals arranged to wipe contaminants from the contacts as first a mechanical, and then an electrical, connection is made between the connecting parts. Another common measure is the use of a so called shuttle pin which occupies in an unconnected configuration a position within the female contact module which will subsequently be occupied by the male contact pin when electrical connection is made. In one typical arrangement the plug is generally cylindrical with an outer housing surrounding a generally cylindrical contact module in which is mounted an axially slidable resiliently biased shuttle pin. The receptacle part is also generally cylindrical and houses a cylindrical male connector pin. When the plug is inserted into the receptacle the male contact pin of the latter axially engages the shuttle pin, and as mechanical engagement is continued the male pin axially displaces the shuttle pin though the contact module until electrical connection is made between the male pin and the female contact module. Typically a wiper seal on the plug wipes the male pin as it penetrates the plug.\nSince well-heads are frequently located at great depth, the connector parts need to be pressure balanced. This is usually achieved by filling chambers in the connector parts with a pressurised dielectric oil, and providing one or more expandable bladders or diaphragms to accommodate movement of the oil as mechanical and electrical engagement is made and unmade.\nAs well-head connections become more complex with increasing requirements for monitoring and control equipment, the space available for connectors of the kind described above becomes reduced, and thus the need for more compact connectors increases."}
{"text": "1. Field of the Invention\nThe present invention relates to an intermediation system providing business information of an article or service through a network and an intermediation device used in the system.\n2. Description of the Related Art\nThe development of the Internet makes it possible to acquire various kinds of information via personal computers, cellular phones and PDAs. To give an example, information on various products or services can be obtained. In addition, an auction site on the Internet or a reverse auction site thereon is used, and various kinds of business is performed (for example, refer to the website at http://auctions.yahoo.co.jp/ (Jan. 7, 2008)).\nAn Internet auction is an auction to be held via the Internet. Specifically, an unspecified number of persons price any of exhibited articles on screens of computers, etc. connected to the Internet, and only the highest bidder can obtain the article. A typical Internet auction site is configured as follows.\n(1) Participants (buyers) bid prices for a certain product individually during a predetermined period. The highest bidder receives an e-mail message indicating that he or she has won the product.\n(2) The exhibitor (seller) of the product informs the highest bidder of the remittance and delivery systems.\n(3) The bidder pays money and the exhibitor sends the product in response to the payment.\nNowadays, there are various kinds of auction sites on the Internet, it is advantageous in that many goods can be purchased cheaply with enjoyment, but because there is a pit for business with an unknown partner, attention is needed. As described above, when simple Internet auctions are utilized, there is the case that a trouble about an article follows.\nCurrently, reverse auctions of which system is opposite to that of auctions are also popular. In a reverse auction, a buyer requests a product that the buyer wants and its desired price. In response, sellers offer their products that are suitable for the request of the buyer. The buyer selects the cheapest one among the offered products, and finally purchases it. A reverse auction system of this type is often applied to public works tenders.\nU.S. Pat. No. 5,794,207 discloses a method of transacting goods on an Internet reverse auction. This method is carried out as follows.\nA buyer sends a mediator of the reverse auction a request of a product that the buyer wants, its conditions, and the buyer's credit card account. In response to the request, sellers send the mediator their estimations. The mediator chooses the most suitable one among the estimations, and informs the buyer about the seller for the selected estimation. The buyer purchases a product that this seller offers.\nAnother style of transaction over the Internet is a website trading. The website trading provides a system in which a website mediates the trade of products between buyers and sellers over the Internet. Products or services that the sellers offer are displayed on this website.\nIn such a website trading, a transaction is conducted as follows.\n(1) Buyers first need to register them to the website as members in order to participate in the trade. A person entitled to be the member is limited to an employee for a corporation or a member for an organization who has purchasing authority.\n(2) An administrator of the website exams the entry persons and, then provides IDs and passwords if the examination is successful. The member (buyer) inputs a desired product and its categories, the number of the product and its price to the site. If some of other members (sellers) have interest in the inputted information, then they can contact the buyer. The sellers offer their products, and a seller who has offers the product of the lowest price among the sellers can bid.\nThe merit of such website trading is that sellers can promote their sales, that is, to get a chance to sell their products to different persons from conventional customers without incurring cost. In addition, it is advantageous to sell competitive products or inventories. Meanwhile, the merit of buyers is to increase their opportunities to purchase desired products at low prices.\nAs described above, electronic commerce or transaction thereof over a network such as the Internet has become popular.\nJP2004-171198 (claims 1 and 18) discloses a method of trading real estate properties. This method makes it possible to make a contract of transacting real estate properties over a network while eliminating unnecessary information, and to provide a real estate trading site containing fresh, secure and a lot of information. Specifically, this method relates to real estate trade to be made between browsers of the site and real estate agencies over a network. Furthermore, the method includes the steps of registering a user to a server connected to a network as a formal member, registering an agency to the server as an agent member, retrieving data on all the agent members from a terminal connected to the server, selecting desired one among all the agent members based on the retrieved data, and making a user confirm a contract between the user and the selected agent member.\nJP2004-171198 (claims 1 and 18) discloses a method of determining population ranking of real estate properties over a network by users and agencies. This method includes the steps of retrieving, by a user, information on real estate properties from a terminal connected to a server through a network, viewing details of attracting properties, viewing drawings of the properties, counting the number of reviewing the details and drawings of the properties, retrieving, by an agency, information on real estate properties from a terminal connected to a server through the network, viewing details of attracting properties, reviewing drawings of the properties, counting the number of reviewing the details and drawings of the properties, and showing the counted number to the users and agencies.\nAs described above, various transaction methods which are conducted over a network are presently available. However, has not yet been conceived, any network system in which sellers can find potential customers based on information on products or services that buyers wish to purchase. Especially, it is difficult to give proper information to potential buyers, because some of buyers who are seemingly earnest might be in fact just browsers.\nTypically, when receiving inquiries from many customers, a seller tries to give the answers to them in the order of the degree of their enthusiasms. However, there is no method of measuring the enthusiasms of buyers. Therefore, a seller may end up in losing potential customers, because the seller gives a higher priority to just viewers such as pricer than potential customers. Also, potential customers may feel dissatisfied because they cannot obtain necessary information timely.\nTaking the above disadvantages into account, the present invention has been conceived. An object of the present invention is to provide a system for mediating transaction information, which can measure enthusiasms of buyers and uses the measured enthusiasms as parameters. An additional object of the present invention is to provide a device used in the above system."}
{"text": "1. Field of the Invention\nThis invention relates to flexible Portland cement polymer concrete and methods for making flexible Portland cement polymer concrete.\n2. Description of the Related Art\nConventional concrete is a dense hard hydrated mass produced from mixtures of Portland cement, coarse aggregate, fine aggregate, and water. Concrete mixtures generally use only sufficient water to make the mixture workable for placement and to yield hardened concrete having a compressive strength of at least 8.3 MPa (1200 psi) after 28 days. Portland cement is a well known cement that upon mixing with water binds or unites the other materials present in the mixture into concrete. The Portland cement is typically a Type I, II, III, IV and V Portland cement as defined in ASTM C 150. The coarse aggregate conventionally comprises particles such as gravel, quartzite, granite, limestone, shale, and other minerals having a particle size greater than 9.5 millimeters (0.375 inches). The fine aggregate employed in Portland cement concretes is most often sand comprised of particles less than about 0.375 inches (9.5 millimeters) in size, typically equal to or less than about 0.1875 inches (4.76 millimeters) in size. Typically, fresh concrete has mixing water exceeding the amount needed for hydration for purposes of workability, handling, and finishing.\nThere are many reasons for the popularity of concrete. It is relatively inexpensive, capable of taking on the shape of a mold, has exceptionally high compression strength and is very durable. However, as a building or construction material, concrete, whether it is reinforced or not, has certain disadvantages. One disadvantage is that concrete has relatively low tensile strength and therefore has little ability to flex without cracking.\nA concrete road surface is a good example where concrete, under repeated loads and flexing due to vehicle traffic, eventually degrades, leading to cracks. This illustrates two primary limitations of Portland cement concrete. The first fundamental limitation is a relatively low tensile strength which primarily exhibits itself as cracks in concrete products. This limitation is addressed by adding reinforcement to the concrete typically in the form of steel fibers, steel bars or steel mesh and/or by controlling the cracking through predetermined crack locations. However, concrete structures can require large amounts of reinforcing material, thereby significantly increasing labor and material costs. The second fundamental limitation is that water intrusion into the concrete can cause both physical damage during freezing and thawing and chemical damage by instigating a series of secondary chemical reactions that can cause degradation of the concrete matrix. This type of behavior can cause surface scaling, cracking and crumbling of concrete much earlier than weather and age might be expected to cause similar degradation. Concrete that has a greater tensile strength, and is thereby able to flex more without failing, would result in surfaces and structures that are less susceptible to wear and cracking.\nWhat is needed therefore is a method for making a concrete having greater tensile strength without requiring the addition of reinforcing materials."}
{"text": "The production of polyesters and copolyesters of dicarboxylic acids and aliphatic glycols has been carried out commercially for many decades. Among the earliest disclosures relating to the production of polyesters and copolyesters is the disclosure in U.S. Pat. No. 2,465,319, issued Mar. 22, 1949. Since this disclosure many variations have been made in the process and many catalysts have been discovered and patented. On Dec. 8, 1970, there issued U.S. Pat. No. 3,546,179, which is directed to the use of compounds containing both silicon and phosphorus atoms a compounds as a catalyst for the production of such polyesters and copolyesters.\nThe use of metal halides and certain silicon compounds as catalysts in the polycondensation of dicarboxylic acids and aliphatic glycols is disclosed in U.S. Pat. Nos. 4,143,057; 4,254,241; and Re-30,554. The metal halides employed in the examples therein are necessarily employed with solvent present. In addition these compounds tend to be hydrolytically unstable and, thus, require careful handling during their use. Thus, although these catalysts are highly advantageous it would be highly desirable to have a catalyst that may be of use without a solvent and which is hydrolytically stable."}
{"text": "In the manufacture of microelectronic devices, integrated circuits are formed on a semiconductor substrate, which generally includes silicon or other materials such as gallium arsenide, and a plurality of interconnect layers such as copper layers, aluminum layers, and the like. With the development of semiconductor technology and integrated circuits, semiconductor technology continues to advance to smaller technology nodes, a semiconductor substrate may include an increased number of integrated circuits. The processed semiconductor substrate is then cut into a number of finished integrated circuits that can be individually packaged.\nHowever, with the continuous shrinking of the size of the semiconductor devices, in processing semiconductor devices with technology nodes of 40 nm and below, the finished products after the laser grooving and diamond cutting processes suffer cracking, peeling, and delamination defects.\nTwo methods are currently available to solve delamination problems: the first one is to form an opening in the top layer of the passivation layer in the scribe line region; and the second one is to reduce the size of the test pad and copper metal layer in the scribe line to avoid stress due to laser cutting. However, the conventional methods cannot prevent delamination from occurring.\nFIG. 1 is a cross-sectional view of a semiconductor device containing a sealing ring, as known in the prior art. A semiconductor device (alternatively referred to as integrated circuit or chip) 10 includes semiconductor substrate 100, a chip edge 110, a cutting region disposed at an outer periphery of the semiconductor device, and a seal ring 120 disposed between the cutting region and a chip region to effectively prevent cracking and peeling in the chip caused by dicing of the semiconductor substrate. The seal ring is formed in the metal layer and the dielectric layer of a semiconductor substrate. The seal ring may be a stacked structure of several metal layers connected together through the through-holes (vias) in the interlayer dielectric layers or other mechanical stress buffering structures.\nDelamination occurs because stress is generated when a metal pad is cut during the laser scribing process, the stress is released after transmitting to the chip. Conventional seal rings cannot stop delamination, low dielectric constant materials are used in next-generation semiconductor processes, which greatly reduce the fracture toughness and adhesion between the layers made of different materials, thus expanding the range of the delamination and cracking in the dielectric layers. Delamination and cracking are generally generated at the edge of the chip and extend to the middle of the chip, adversely affecting the functional integrity and yield of the chip.\nThus, there is a need to provide a new semiconductor device that overcomes the disadvantages associated with prior art."}
{"text": "The present invention relates to centrifugal pumps in general, and more particularly to improvements in means for reducing or eliminating the effects of cavitation in part-load region of operation of centrifugal pumps. Still more particularly, the invention relates to improvements in means for reducing or eliminating cavitation in centrifugal pumps of the type wherein the shaft which drives the impeller or impellers extends through the fluid-admitting inlet of the pump.\nIn a conventional centrifugal pump which is to operate within a wide range of capacities, and wherein the optimum rate of inflow of fluid takes place at the design duty point, the parts undergo extensive wear owing to cavitation whenever the pump is operated in part-load region. The undesirable effects of cavitation are more pronounced in larger centrifugal pumps, and the wear owing to cavitation increases with increasing rotational speed. Attempts to reduce the wear which is attributable to cavitation include limiting the interval of operation in the part-load region, utilizing highly wear-resistant materials for those parts which are most likely to be affected by cavitation and/or providing inserts which are installed in regions where the effects of cavitation are most likely to induce rapid wear or destruction and to provide for replacement of such inserts at regular intervals or when the need arises.\nThe reasons for pronounced cavitation during operation in the part-load region are manifold. They include the following: When the rate of flow (capacity) is less than 60-70 percent of the flow rate at the design duty point, a ring-shaped eddy develops in the outer zone of the intake of the impeller (such eddy is known as part load eddy). The flow of fluid which forms the eddy is in a direction away from the impeller and counter to the direction of inflowing fluid. The intensity of the eddy increases in response to decreasing rate of flow into the impeller, i.e., the eddy constricts the inlet at a rate which increases with decreasing capacity. The annular stream of inflowing fluid, whose outer radius decreases with decreasing rate of flow, enters the pump at the same speed as when the operating point matches the design duty point; however, the relative speed of fluid in the hub region of the impeller is lower.\nThe just discussed phenomenon entails a reduction of the NPSH (net positive suction head) for incipient cavitation (NPSH.sub.i) when the pump is operated in the part-load region, i.e., when the operating point is such as to permit for development of part load eddies. The formation of a part load eddy is affected by the presence of ribs, elbows or other structural elements at the inlet, and such influence upon the eddy results in an increase (rather than a decrease) of NPSH.sub.i to a multiple of the anticipated value. Thus, damage owing to cavitation is very pronounced especially since (and as already pointed out above) the optimum rate of fluid flow takes place only at the design duty point."}
{"text": "Online social network applications are becoming increasingly popular. People rely on such applications to keep in touch and interact with friends and family. Typically, to utilize such an online social networking application, an individual that logs into a web site that corresponds to the application will be provided a visual depiction of messages sent to the individual as well as status updates of contacts (friends) of the individual.\nWhile these online social networking applications are useful tools that allow people to interact with their contacts, often times messages transmitted using social networking applications are relatively short with little context associated with them. Thus, an assumption exists that an individual receiving a relatively short message will have some knowledge of what the sender of the message is referring to in the message."}
{"text": "The Word Wide Web enables a web resource, such as a document, to be identified by a web address, such as a Uniform Resource Locator (URL) that provides a global naming scheme for web resources. A web address specifies a web resource's location on the World Wide Web and explains the means by which a web resource can be accessed via the Internet. For example, a URL for a web page may utilize Hypertext Transfer Protocol (HTTP), and when the URL is entered into an address bar of a web browser, then the web browser displays the web page. The web page may include information including a hypertext link, such as Hypertext Markup Language (HTML), that provides for the web page to link to other web pages. It is these hyperlinks of web pages that create what can be called a web of information, which is another description of the World Wide Web.\nAlthough it is not possible to know how many web sites there are on the World Wide Web, an estimate is that the count is above one billion. Such reflects an immense amount of information for a user to manually browse with a web browser, such as Microsoft Internet Explorer, Google Chrome, and the like. In order to assist users in finding and processing information available on the World Wide Web, a process known as web crawling has been developed to automate the browsing, grabbing and sifting of information on the World Wide Web. Such can attempt to simulate a user's exploration of the World Wide Web to automatically extract useful information. Web crawling is accomplished with a web robot that is a software application of process instructions that execute automated and repetitive tasks via the Internet. One of those repetitive tasks is to extract structured and unstructured information from a web page, and convert that information into a form that can be stored in a central database. Such data extraction from a web page is commonly called web scraping, and may be beneficial for online price comparisons, centralized job postings on career web sites, and analyzing market research, among other uses. Web scraping allows for the collection of information from a web site without having to directly access any back-end databases, for example, through applicable application programming interfaces (APIs). Traditional financial management systems provide a mechanism to upload information from a service provider's system to the financial management system, such as through a data transfer interface that may utilize APIs. For example, a banking system may transfer a user's bank transactional records to a financial management system via APIs associated with the Open Financial Exchange (OFX) Standard or the like.\nWeb sites are generally intended to be displayed to an end-user and are not generally intended to be utilized for data transfers between software systems. Nevertheless, scraping is a data transfer process that can be at a semantic level or a syntactic level. Semantic level scraping utilizes a common framework for structured data. A web page may utilize a markup language such as Extensible Markup Language (XML) that promotes structured data with rules for encoding data that is both human-readable and machine-readable. A script may be created from a scripting language such as PERL, Tcl, Python, and the like, and a script may extract well-formed XML content comprising data at the semantic level. Under some embodiments, a script comprises processing instructions. Such may be advantageous because a web page's HTML may change due to changes in the look and feel of the web page, yet the script can still extract the XML content at a semantic level. For example, the data may be presented in a systematic way through XML elements that convey the meaning of the data.\nOn the other hand, syntactic level scraping refers to scraping data that is not presented in a systematic way. Web scraping at the syntactic level is often considered a manually intensive process to create scripts to extract information. For example, one may rely on knowledge of specific markup practices or document structures to determine data items among other web page elements. Such knowledge can be used to recognize data structures on a web page, and create a script to extract data at the syntactic level. Scraping data that is presented in a syntactic way can be prone to difficulties because when a web page's layout or HTML changes, the script must often be rewritten in response to continue to scrape the web page at the syntactic level. Until the script is rewritten, the data that had been retrieved in the past can often no longer be retrieved.\nThe creation of a script to scrape a web page at the syntactical level may be further complicated by the lack of domain-specific knowledge for a web page, in which a domain may be an industry grouping such as financial institutions, healthcare institutions, insurance institutions, travel institutions, legal institutions, and others. Having knowledge to create scraping scripts in one domain does not necessarily mean that one can create scraping scripts in other domains. For example, the creation of a script to scrape a financial institution web page at the syntactic level may require knowledge about how financial transactions are presented on a web page. In contrast, the creation of a script to scrape a healthcare institution web page at the syntactic level may require knowledge about how healthcare transactions are presented on a web page. Because of the diversity of domains, a developer who creates scripts and who possesses the knowledge to extract transactions from one domain may not have the knowledge to extract transactions from another domain.\nAccordingly, there is a need to improve the creation of scraping scripts that extract information from web pages across different domains."}
{"text": "Large commercial irrigation systems such as those used on golf courses or croplands use sprinklers, sensors or other components that are normally powered from 24 V AC power lines that can be several miles long and can serve many hundreds of components. Various systems have been proposed for powering and controlling the components of such a system with just two wires. For example, U.S. Pat. No. 3,521,130 to Davis et al., U.S. Pat. No. 3,723,827 to Griswold et al., and U.S. Pat. No. 4,241,375 to Ruggles disclose systems in which sprinkler valves along a cable are turned on in sequence by momentarily interrupting the power or transmitting an advance signal from time to time.\nA problem with this approach is that it does not allow the operator to freely turn on or off any selected sprinkler or set of sprinklers at different times. This problem is typically resolved by providing separate controllers in the field to operate groups of sprinklers in accordance with a program stored in them, or transmitted to them by radio or other means. Alternatively, it has been proposed, as for example in U.S. Pat. No. 3,578,245 to Brock, to operate individual sprinkler sets from a central location by superimposing a frequency-modulated signal or DC pulses onto the 24 V AC power line.\nAll of these approaches are expensive. For example, a system with hundreds of sprinklers requires miles of expensive, heavy wiring to accommodate the current drawn by a large number of valves that may be watering simultaneously. Also, heavy use of D.C. current may cause electrolysis issues with electrical components.\nOne alternative to these traditional irrigation systems are two-wire power and communications systems, such as the system shown in U.S. Pat. No. 7,358,626, the contents of which are incorporated by reference. In such systems, two wires supply both A.C. power and digital control communications from a controller to a plurality of decoders.\nWhile these A.C. power and digital communication systems generally work well, they also have several disadvantages. First, decoder circuitry that listens for communications are sensitive to power surges by lightning. Hence, expensive, external surge devices with ground rods must be installed at short intervals along the power and communication wires.\nSecond, some two wire power and communication systems use bursts of current by valve or sensor decoders to acknowledge or otherwise communicate with a central controller. These decoders relied on fixed resistor circuits to switch to a circuit path with a resistor between the common current and ground on the decoder. Since these resistor circuits resist at a fixed value, any reduction in the voltage causes a proportional reduction in the current pulse (e.g., due to Ohm's Law). Therefore if the current burst of the acknowledgement pulse is reduced sufficiently, the gateway 16 will not recognize it and therefore will falsely determine that the decoder did not receive a command. As decoders are installed further away from a gateway or central controller, the resistance of the power and communication wires becomes significant and reduces the voltage available across the decoder. Additionally, large numbers of decoders may also further reduce voltage and increase resistance on the power and communication wires. Hence, these systems effectively become limited in size and in the number of decoders."}
{"text": "Various different types of tire tread wear indicators are known in the prior art to indicate when a tread has worn to a level where the tire is no longer safe for use. On such tread wear indicator is shown in U.S. Pat. No. 5,303,756 issued to J L Hill which shows the use of wear indicating members in the shape of pins or pegs molded in the tread material adjacent the inner surface of the tread. These members are made in a color contrasting to the color of the tread so that as the tread wears to a certain level the wear indicators become visible. U.S. Pat. No. 4,226,274 issued to H Y Awaya, et al. shows another example of the use of contrasting colored indicia embedded in the tread to indicate the amount of tread remaining and when the tread has worn to an unsafe level. Likewise U.S. Pat. No. 3,516,467 issued to C D Simms shows a tire tread having annular ribs of a contrasting colored material embedded at the radially inner portion of the tread in such a position that when the tread wears to a certain depth, the ribs become visible. A similar concept is shown in U.S. Pat. No. 3,833,040 issued to C G Bins. U.S. Statutory Invention Registration No. H1283 shows a tire having plies of contrasting colors to each other and with a bottom ply having a contrasting pattern so that various levels of wear are indicated by the different colored plies and the bottom ply serves as a warning of an unsafe condition. U.S. Pat. No. 2,706,509 issued to E J White indicates various levels of tire wear by means of stepped surfaces located in the voids or grooves between adjacent ribs of a tire. This indicator will show for example of whether the tread is worn one-forth, one-half or to some other level for the purpose of providing a price adjustment on a tire which may fail before the tread is fully worn.\nWhile the above described tread wear indicators may function as various levels of effectiveness for the purposes for which they were designed, none of the above wear indicators was designed for the purpose of indicating when a tire should be rotated from one wheel location to another on a vehicle to compensate for any uneven tread wear and extend the life of the tire."}
{"text": "1. Filed of the Invention\nThe present invention relates to a transfer apparatus providing and a molded article production apparatus having the transfer apparatus. The transfer apparatus provides a component for insert molding to a die by insert molding, and removes a molded article which is formed by the die from the die.\n2. Description of the Related Art\nGenerally, various electronic devices, for example a light such as headlights and a rear light, and a motor such as a starter motor and a motor for air condition, are mounted in a vehicle as a mobile object.\nA junction block is arranged in appropriate place of the vehicle so as to provide electrical power into the above mentioned various electronic devices. The junction block is included by putting together various electronic circuit units of multiple fuses and relays.\nA junction block has a fuse, a relay or a bus bar. Thereby, the junction block is called a fuse block, a relay box, or an electronic connection box. In the present invention, the above fuse block, the relay box and the junction block are called an electronic connection box as a generic name in the below.\nThe above mentioned electric connection box has a box body forming an outer case of it and a wiring block. The box body is made of an insulation synthetic resin, and is formed into a box shape. Electrical part, for example a relay and a fuse, and a connector of a wiring harness are mounted on the box body.\nThe wiring block has a conductive bus bar, which is buried into the box body, and a terminal. The wiring block connects electrically each other the electrical part, which are attached to the box body, and the terminal of the connector of the wiring harness in a predetermined pattern.\nThe bus bar and the terminal constituting the wiring block are set into a cavity of die for insert molding, and the cavity is filled with synthetic resin. Thereby, the bus bar and the terminal are integrally formed with the box body together. That is, the bus bar and the terminal are integrally formed with the box body by insert molding.\nHeretofore, the above mentioned insert molding is formed by various molding production apparatus (for example, see a patent document 1). A molded article production apparatus shown in the patent document 1 has an insert molding machine, a robot as a transfer apparatus providing components into a die and a robot as a transfer apparatus removing a formed molding from the die.    Patent document 1:            Japan published patent application H9-201849        "}
{"text": "Communication systems are known to support wireless and wire lined communications between wireless and/or wire lined communication devices. Such communication systems range from national and/or international cellular telephone systems to the Internet to point-to-point in-home wireless networks. Each type of communication system is constructed, and hence operates, in accordance with one or more communication standards. For instance, wireless communication systems may operate in accordance with one or more standards including, but not limited to, IEEE 802.11, Bluetooth, advanced mobile phone services (AMPS), digital AMPS, global system for mobile communications (GSM), code division multiple access (CDMA), local multi-point distribution systems (LMDS), multi-channel-multi-point distribution systems (MMDS), radio frequency identification (RFID), Enhanced Data rates for GSM Evolution (EDGE), General Packet Radio Service (GPRS), and/or variations thereof.\nDepending on the type of wireless communication system, a wireless communication device, such as a cellular telephone, two-way radio, personal digital assistant (PDA), personal computer (PC), laptop computer, home entertainment equipment, RFID reader, RFID tag, et cetera communicates directly or indirectly with other wireless communication devices. For direct communications (also known as point-to-point communications), the participating wireless communication devices tune their receivers and transmitters to the same channel or channels (e.g., one of the plurality of radio frequency (RF) carriers of the wireless communication system or a particular RF frequency for some systems) and communicate over that channel(s). For indirect wireless communications, each wireless communication device communicates directly with an associated base station (e.g., for cellular services) and/or an associated access point (e.g., for an in-home or in-building wireless network) via an assigned channel. To complete a communication connection between the wireless communication devices, the associated base stations and/or associated access points communicate with each other directly, via a system controller, via the public switch telephone network, via the Internet, and/or via some other wide area network.\nFor each wireless communication device to participate in wireless communications, it includes a built-in radio transceiver (i.e., receiver and transmitter) or is coupled to an associated radio transceiver (e.g., a station for in-home and/or in-building wireless communication networks, RF modem, etc.). As is known, the receiver is coupled to an antenna and includes a low noise amplifier, one or more intermediate frequency stages, a filtering stage, and a data recovery stage. The low noise amplifier receives inbound RF signals via the antenna and amplifies then. The one or more intermediate frequency stages mix the amplified RF signals with one or more local oscillations to convert the amplified RF signal into baseband signals or intermediate frequency (IF) signals. The filtering stage filters the baseband signals or the IF signals to attenuate unwanted out of band signals to produce filtered signals. The data recovery stage recovers raw data from the filtered signals in accordance with the particular wireless communication standard.\nAs is also known, the transmitter includes a data modulation stage, one or more intermediate frequency stages, and a power amplifier. The data modulation stage converts raw data into baseband signals in accordance with a particular wireless communication standard. The one or more intermediate frequency stages mix the baseband signals with one or more local oscillations to produce RF signals. The power amplifier amplifies the RF signals prior to transmission via an antenna.\nWhile transmitters generally include a data modulation stage, one or more IF stages, and a power amplifier, the particular implementation of these elements is dependent upon the data modulation scheme of the standard being supported by the transceiver. For example, if the baseband modulation scheme is Gaussian Minimum Shift Keying (GMSK), the data modulation stage functions to convert digital words into quadrature modulation symbols, which have constant amplitude and varying phases. The IF stage includes a phase locked loop (PLL) that generates an oscillation at a desired RF frequency, which is modulated based on the varying phases produced by the data modulation stage. The phase modulated RF signal is then amplified by the power amplifier in accordance with a transmit power level setting to produce a phase modulated RF signal.\nAs another example, if the data modulation scheme is 8-PSK (phase shift keying), the data modulation stage functions to convert digital words into symbols having varying amplitudes and varying phases. The IF stage includes a phase locked loop (PLL) that generates an oscillation at a desired RF frequency, which is modulated based on the varying phases produced by the data modulation stage. The phase modulated RF signal is then amplified by the power amplifier in accordance with the varying amplitudes to produce a phase and amplitude modulated RF signal.\nAs the desire for wireless communication devices to support multiple standards continues, recent trends include the desire to integrate more functions on to a single chip. However, numerous functions increase the amount of power consumed by the IC. Thus, better methods of managing power consumption and system resources within a multiple function IC are desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit using an oxide semiconductor and a method of driving the semiconductor integrated circuit.\n2. Description of the Related Art\nIn a semiconductor integrated circuit, a NAND circuit, a NOR circuit, an inverter circuit (NOT circuit), and the like are used in combination.\nA conventional logic circuit includes a CMOS circuit in which p-channel transistors and n-channel transistors are used in a complementary combination. The conventional logic circuit loses data when power supply is stopped.\nIn addition, Patent Document 1 discloses a logic circuit which is formed with a CMOS circuit configuration using a spin MOSFET to maintain a nonvolatile logic function. As one example thereof, a circuit diagram of an XOR circuit from which a result of exclusive OR of one input and another input is obtained is disclosed in Patent Document 1."}
{"text": "The present invention relates to a process for preparing conveniently an ascorbic acid-2-monophosphate salt, particularly an L-ascorbic acid-2-monophosphate (hereinafter sometimes abbreviated to xe2x80x9cAsMPxe2x80x9d) salt in a high yield, by hydrolyzing an ascorbic acid-2-polyphosphate or a salt thereof, particularly an L-ascorbic acid-2-polyphosphate (hereinafter sometimes abbreviated to xe2x80x9cAsPPxe2x80x9d) or a salt thereof, as the starting material in the presence of magnesium ion.\nL-Ascorbic acid (vitamin C) is known to have various physiological and pharmacological actions, and particularly by virtue of its effect of preventing melamine pigmentation, L-ascorbic acid has been used in whitening cosmetics. However, L-ascorbic acid is unstable to oxygen or heat, therefore, L-ascorbic acid derivatives stable against oxygen or heat have been heretofore formed by converting the hydroxyl group at the 2-position into a phosphate.\nA salt, particularly a magnesium salt of L-ascorbic acid-2-monophosphate (hereinafter the xe2x80x9cmagnesium L-ascorbic acid-2-monophosphatexe2x80x9d is sometimes abbreviated to xe2x80x9cAPMxe2x80x9d) is being used as a stabilized vitamin C derivative. APM exhibits excellent stability in cosmetic materials and scarcely decomposes. Furthermore, APM is easily absorbed through the skin and by the action of phosphatase present inside the human body, L-ascorbic acid is liberated to bring about various physiological actions such as prevention of melamine pigmentation.\nHigh-purity AsMP salts are useful as a stabilized derivative of L-ascorbic acid and can be used in cosmetics, medical products, food additives, feed and other various industrial fields.\nFor monophosphorylating L-ascorbic acid, three methods described below are known. A first method comprises monophosphorylating L-ascorbic acid using phosphorus oxychloride as the phosphorylating agent and this method is described, for example, in JP-B-45-30328 (the term xe2x80x9cJP-Bxe2x80x9d as used herein means an xe2x80x9cexamined Japanese patent publicationxe2x80x9d), JP-B-52-18191, JP-B-59-4438 and JP-A-2-27969 (the term xe2x80x9cJP-Axe2x80x9d as used herein means an xe2x80x9cunexamined published Japanese patent applicationxe2x80x9d) and the like. The second method comprises transposing the phosphoric acid group of a phosphoric acid donor (e.g., adenosine triphosphate, pyrophosphoric acid and the like) to L-ascorbic acid using an enzyme and this method is described, for example, in JP-A-2-283283. The third method comprises reacting a soluble salt of metaphosphoric acid used as a phosphorylating agent with L-ascorbic acid and drying the aqueous solution to obtain AsMP salt and this method is described, for example, in JP-A-5-155893.\nThe first method above is most frequently used and according to Carbohydrate. Res., 67, 127-138 (1978), L-ascorbic acid-2-monophosphate is obtained as a main product in the form of tricyclohexylamine salt in a yield of 86%. The by-products are L-ascorbic acid-3-monophosphate, L-ascorbic acid-2-diphosphate (hereinafter sometimes abbreviated to xe2x80x9cAsDPxe2x80x9d) and 2,2xe2x80x2-bis(L-ascorbic acid)phosphate. For the purification thereof, a complicated step such as ion exchange chromatography is necessary. L-ascorbic acid has four hydroxyl groups active in the reaction with phosphorus oxychloride and phosphorus oxychloride has three active sites, therefore, many by-products are produced. Furthermore, chloride ion generated from phosphorus oxychloride after the reaction is 3 molar times the phosphorus oxychloride used and this requires purification by electrodialysis and the like.\nThe second method is advantageous from the viewpoint that generation of by-products is prevented but still has the defect that the productivity is low. According to JP-A-2-283283, AsMP is produced using an enzyme fixed to an ion exchange resin and using diphosphoric acid (pyrophosphoric acid) as a phosphoric acid donor. However, the yield of AsMP of 34% is very low. Also, the amount of AsMP produced is as low as 1.7% in terms of the concentration in the reaction solution and large scale reaction equipment is necessary. Furthermore, in order to obtain high-purity AsMP salt, purification through many steps such as removal of a large amount of water and recovery of unreacted L-ascorbic acid is required.\nThe third method is a method of drying L-ascorbic acid-2-triphosphate (hereinafter sometimes abbreviated to xe2x80x9cAsTPxe2x80x9d) in an aqueous solution and thereby stepwise hydrolyzing the phosphate bonded to produce AsMP. According to JP-A-5-155893, AsMP is produced by preparing an aqueous solution of AsTP from a soluble salt of metaphosphoric acid and L-ascorbic acid, adjusting the pH thereof to about 5.5 to 6.5 and performing drying at a temperature of from 120 to 180xc2x0 C. As known from the pH of aqueous solution, this stepwise hydrolysis reaction of phosphate is acid hydrolysis reaction. Therefore, although it is set forth therein that sodium and calcium are advantageous as the cation present in the aqueous solution, the alkali property of the cation has no relation to the acid hydrolysis, since the pH of the aqueous solution is acidic.\nFurthermore, the product AsMP obtained by this production method has extremely low purity and the yield is also low. As seen from the composition of the product obtained by this method described in the Examples, the AsMP purity is 61% and the yield based on AsTP is 66%. In addition, much L-ascorbic acid, AsDP and AsTP are present as by-products. The presence of L-ascorbic acid shows that the acid hydrolysis reaction of AsTP extended to L-ascorbic acid over AsMP.\nThis means that it is useless to phosphorylate the hydroxyl group at the 2-position of L-ascorbic acid so as to improve lack of stability of L-ascorbic acid to oxygen or heat. Furthermore, the L-ascorbic acid present in the product decomposes due to oxygen or heat causing coloration or produce oxalic acid which is a skin irritant. The product cannot be used in cosmetics and the like as it is. Therefore, the impurities must be purified by a complicated method such as ion exchange chromatography. For these reasons, this production method of AsMP by the drying of AsTP aqueous solution is not an industrial attractive.\nAs described above, AsMP salt is useful as a stabilized derivative of L-ascorbic acid, nevertheless, the production methods thereof are complicated and disadvantageous in industry. Under these circumstances, a simplified preparation process is demanded therefor.\nAs a result of extensive studies to eliminate the above-described defects, the present inventors have found that high-purity AsMP salt can be conveniently prepared in a high yield by hydrolyzing, preferably under heating AsPP, or a salt thereof at a pH of 7 or more as a starting material in the presence of magnesium ion. The present invention has been accomplished based on this finding of a new process.\nThe present invention provides the following embodiments for preparing an ascorbic acid-2-monophosphate salt:\n(1) a process for preparing an ascorbic acid-2-monophosphate salt, comprising hydrolyzing an ascorbic acid-2-polyphosphate or a salt thereof in the presence of magnesium ion at a pH of 7 or more;\n(2) the process as described in (1) above, comprising hydrolyzing a reaction product of an ascorbic acid and a polyphosphoric acid or a salt thereof in the presence of magnesium ion at a pH of 7 or more;\n(3) the process as described in (1) above, comprising eating a mixture comprising:\n(A) an ascorbic acid-2-polyphosphate or a salt thereof,\n(B) water,\n(C) a magnesium compound capable of supplying magnesium ion in water, and\n(D) a base necessary for maintaining a pH of 7 or more;\n(4) the process as described in (3) above, wherein the\nascorbic acid-2-polyphosphate or a salt thereof (A) is one or more selected from the group consisting of an L-ascorbic acid-2-triphosphate and a salt thereof, an L-ascorbic acid-2-diphosphate and a salt thereof, and a mixture thereof;\n(5) the process as described in (3) above, wherein the magnesium compound (C) is one or more selected from the group consisting of magnesium chloride, magnesium sulfate, magnesium acetate, magnesium nitrate, magnesium oxide, magnesium hydroxide, magnesium carbonate and a mixture thereof;\n(6) the process as described in (3) above, wherein the base (D) is one or more selected from the group consisting of an alkali metal hydroxide, an alkali metal carbonate, an alkaline earth metal hydroxide, an alkaline earth metal oxide, an alkaline earth metal carbonate, a tertiary amine and a mixture thereof;\n(7) the process as described in (3) above, wherein the compound serving both as the magnesium component (C) and the base (D) is one or more selected from the group consisting of magnesium hydroxide, magnesium oxide and magnesium carbonate;\n(8) the process as described in (3) above, wherein the heating temperature is from 35xc2x0 C. to the boiling point of the aqueous reaction solution;\n(9) the process as described in (3) above, which includes, after the completion of reaction, adjusting the pH by adding an acid to the reaction solution or the diluted or concentrated solution to precipitate a phosphate salt and removing the phosphate salt;\n(10) the process as described in (3) above, which includes, after the completion of reaction, adjusting the pH by adding an acid to the reaction solution or the diluted or concentrated solution to precipitate a phosphate salt and dissolving the precipitated phosphate salt and adding a base to the dissolved solution to precipitate the phosphate salt and removing the phosphate salt;\n(11) the process as described in (9) above, which includes crystallizing an ascorbic acid-2-monophosphate salt from the aqueous reaction solution or an aqueous solution after the removal of phosphate using a water-soluble organic solvent; and\n(12) the process as described in (11) above, which includes drying the wet ascorbic acid-2-monophosphate salt crystallized from the aqueous reaction solution or an aqueous solution after the removal of phosphate using a water-soluble organic solvent to obtain the ascorbic acid-2-monophosphate salt as a powder.\nIn the present invention, the ascorbic acid structure may be any of an L-form, a D-form and a racemic compound, however, the present invention is described below by referring to the L-form which is industrially useful.\nAccording to the process of the present invention, an aqueous solution containing AsPP or a salt thereof is hydrolyzed, preferably under heating, at a pH of 7 or more in the presence of magnesium ion and thereby the phosphoric acid ester bond of AsPP is stepwise alkali-hydrolyzed to prepare AsMP.\nThis preparation process of the present invention comprises allowing magnesium ion to be present together and adjusting the pH to 7 or more, the stepwise hydrolysis from AsPP to AsMP proceeds almost constantly in quantity and the hydrolysis reaction from AsMP to L-ascorbic acid scarcely proceeds.\nIn the preparation process of the present invention, the ion allowed to be present at the hydrolysis cannot be just any alkaline earth metal ion but must be magnesium ion. The magnesium ion has an action of accelerating the hydrolysis reaction from AsPP to AsMP but not accelerating the hydrolysis from AsMP to L-ascorbic acid. The reasons therefor are not clearly understood but the magnesium ion maintains relatively mild alkalinity even in a strong alkali condition and has low water-solubility and these properties are considered to take a part in the above-described action.\nFor example, when magnesium oxide was added in 1.7 molar times to an aqueous solution containing AsMP and the solution was heated at 60xc2x0 C. for 6 hours, the amount of residual AsMP was 100 wt %. When the same operation was performed using calcium oxide in place of magnesium oxide, the amount of residual AsMP was 41 wt % and this clearly showed formation of L-ascorbic acid and furthermore formation of oxalic acid by the decomposition of L-ascorbic acid. In the examination where magnesium chloride was added in 1.7 molar times to a solution containing AsMP sodium salt and heated at 60xc2x0 C. for 6 hours, the amount of residual AsMP was also 100 wt %. However, when the same operation was performed using calcium chloride in place of magnesium chloride, the amount of residual AsMP was 63 wt %. The same results were obtained in the examination using barium ion in place of calcium ion.\nAs such, magnesium ion has an action of not accelerating the hydrolysis of AsMP but calcium ion and the like have no such an action. which can be used in the present invention include polyphosphoric acids such as triphosphoric acid, diphosphoric acid (pyrophosphoric acid) and metaphosphoric acid, and salts thereof. Accordingly, the AsPP usually includes AsTP, AsDP, and salts and mixtures thereof. These may be prepared by the method described in U.S. Pat. 5,149,829 and the like. In the present invention, the dehydration product of an ascorbic acid and a polyphosphoric acid or a salt thereof may be used as it is in the hydrolysis.\nAccording to U.S. Pat. 5,149,829, an L-ascorbic acid and a polyphosphate salt such as sodium trimetaphosphate were reacted at a temperature of from 33 to 35xc2x0 C. while adjusting the pH to 10.5 to 10.7 by NaOH or KOH. After the completion of reaction, the inorganic phosphoric acid compound was removed by precipitation as a calcium salt using calcium chloride at a temperature as low as 10xc2x0 C. Thereafter, the residue was concentrated, crystallized from ethanol and dried to obtain the objective AsPP salt as a mixture of AsMP salt, AsDP salt and AsTP salt. In the method of U.S. Pat. 5,149,829, the inorganic phosphoric acid compound precipitated is removed, however, in the present invention, inorganic phosphoric acids may be mixed in the AsPP salt. In the case where use of AsPP as an acid but not as a salt is intended, the acid may be obtained by dissolving an AsPP salt in water and decationizing it by passing the solution through a strongly acidic cation exchange resin at a low temperature. Also, when use of only pure AsTP or only pure AsDP as a raw material is intended, the raw material may be obtained by purifying the AsPP using chromatography.\nIn the present invention, the AsPP concentration is, for example, in water, suitably from about 1 to 80 wt %. If the AsPP concentration is low, the reaction slowly proceeds, whereas it is excessively high, the viscosity increases or magnesium salt of AsPP precipitates. Accordingly, the AsPP concentration is preferably from about 5 to 50 wt %, more preferably from about 10 to 30 wt %.\nThereafter, magnesium ion is added to the aqueous solution of AsPP or a salt thereof. In the case where the AsPP is a salt, a water-soluble magnesium compound such as magnesium chloride is preferably used and in the case where the AsPP is an acid, a magnesium compound showing alkalinity such as magnesium oxide is preferably used. Examples of the magnesium ion-supplying compounds include magnesium chloride, magnesium sulfate, magnesium acetate, magnesium nitrate, magnesium oxide, magnesium hydroxide, magnesium carbonate and the like. These compounds may be used individually or as a mixture without any problem.\nThe amount of magnesium ion added is an equimolar amount or more to the phosphoric acid group of AsPP. As the alkali hydrolysis of AsPP proceeds, phosphoric acid is produced in the aqueous reaction solution and the phosphoric acid produced reacts with magnesium ion present in the aqueous reaction solution to form magnesium hydrogenphosphate or magnesium phosphate. These magnesium salts of phosphoric acid have very low solubility in water and accordingly precipitate out. As a result, magnesium ion in the aqueous solution necessary for the hydrolysis reaction decreases. Accordingly, in order to maintain the magnesium ion in the reaction system even when the hydrolysis reaction of AsPP proceeds, the magnesium ion is required in an equimolar amount or more to the phosphoric acid group of AsPP. In this reaction, the amount of magnesium ion added has no upper limit in particular, however, the presence of excess magnesium ion has no effect on the reaction and this is of no benefit. Accordingly, the amount of magnesium added is preferably from about 1.1 to 2.0 mol, more preferably from 1.3 to 1.8 mol, based on the phosphoric acid group of AsPP.\nIn the preparation process of the present invention, the pH is suitably set to be 7 or more. If the pH is less than 7, the stepwise hydrolysis reaction in the present invention becomes hydrolysis accelerated by an acid, as a result, the hydrolysis of the objective product AsMP is also accelerated to decompose it into L-ascorbic acid and phosphoric acid. Accordingly, the aqueous reaction solution must be set to a pH of 7 or more using an appropriate alkali. Examples of the alkali used here include alkali metal hydroxides, alkali metal carbonates, alkaline earth metal hydroxides, alkaline earth metal oxides, alkaline earth metal carbonates, tertiary amines, and mixtures of two or more thereof. In the case where the alkaline earth metal compound used is magnesium oxide, magnesium hydroxide or magnesium carbonate, magnesium ion can be supplied to the aqueous reaction solution and thus actions of adjusting the pH and at the same time accelerating the hydrolysis of AsPP can be achieved. In the present invention, substantially no problem arises as long as the pH is 7 or more. However, the reaction slowly proceeds in the alkali region closer to the neutral, whereas reaction in the strong alkali condition requires alkali-resistant equipment. Accordingly, the pH is preferably from 8 to 13, more preferably from 9 to 11.\nIn addition, heating is preferably performed in the present invention. If the reaction is performed at a low temperature, a very long time is necessary until the completion of reaction and the plant is not utilized efficiently and this is disadvantageous. Accordingly, the reaction is preferably performed at a temperature of from 35xc2x0 C. to the boiling point of aqueous reaction solution, more preferably from 50 to 90xc2x0 C.\nThe thus-obtained AsMP-containing aqueous solution is diluted or concentrated to an appropriate AsMP concentration and a phosphate salt precipitated (example: phosphoric acid alkali salt) is removed by the solid-liquid separation. At this time, the pH is adjusted with an appropriate acid to effectively precipitate the phosphate salt. Examples of said acid are concentrated or diluted hydrochloric acid. Or, the phosphate salt precipitated may be dissolved by an acid and again precipitated by adding a base.\nIn the aqueous solution containing high-purity AsMP salt after the removal of the phosphate salt precipitated, AsMP salt is precipitated using a water-soluble organic solvent. Examples of the water soluble organic solvent used here include methanol, ethanol, isopropanol, tetrahydrofuran, dioxane, acetone and the like. Among these, ethanol is preferred in view of the cost and ease of recovery.\nThe thus-obtained wet AsMP salt is dried by vacuum drying or fluidized bed drying to obtain AsMP salt in a high yield.\nIn particular, AsMP magnesium salt has excellent stability against heat or light, accordingly, it is used not only for cosmetic powders or lotions but also, for example, for medical products (e.g., oral cavity preparation, eyedrops, products for bath and the like), cosmetics (e.g., lotion, emulsion, cream, pack and the like), food (e.g., bread and the like) and animal feed (e.g., feed for breeding prawn, salmon, yellow tail, eel, carp and the like)."}
{"text": "1. Technical Field\nThe present invention relates to a MEMS oscillator including a MEMS vibrator and a method of manufacturing thereof.\n2. Related Art\nFIG. 13 shows the circuit configuration of a related-art MEMS vibrator. One of electrodes of the MEMS vibrator 201 is connected to a coupling capacitance 203a, and the coupling capacitance 203a is connected to an input terminal. The other electrode of the MEMS vibrator 201 is connected to a coupling capacitance 203b, and the coupling capacitance 203b is connected to an output terminal 220. A bias voltage applying circuit 202 is connected between the electrodes of the MEMS vibrator 201 via bias resistances 204a and 204b. \nWhen the MEMS vibrator 201 is operated, a DC bias voltage is applied between the electrodes of the MEMS vibrator 201 by the bias voltage applying circuit 202 including a reference voltage generating circuit so that the MEMS vibrator 201 is free from the influence of an external DC voltage. The input terminal 210 and the output terminal 220 which are coupled to an external circuit and the electrodes of the MEMS vibrator 201 are connected via the coupling capacitances 203a and 203b. To obtain AC insulation between the bias voltage applying circuit 202 which applies the DC bias voltage and the electrodes of the MEMS vibrator 201, the bias resistances 204a and 204b are inserted between the electrodes of the MEMS vibrator 201 and the bias voltage applying circuit 202.\nFIG. 14 shows an exemplary MEMS oscillator including the related-art MEMS vibrator and an oscillator circuit. One of the electrodes of the MEMS vibrator 201 is connected to the coupling capacitance 203a, while the other electrode of the MEMS vibrator 201 is connected to the coupling capacitance 203b. The bias circuit 202 is connected between the electrodes of the MEMS vibrator 201, and a bias voltage input terminal is connected to the bias circuit 202. The coupling capacitance 203a is connected to an input side of an amplifier circuit 230, while the coupling capacitance 203b is connected to an output side of the amplifier circuit 230. The output side of the amplifier circuit 230 and the coupling capacitance 203b are connected to an input side of a buffer circuit (buffer amplifier) 240, while an output side of the buffer circuit 240 is connected to an output terminal.\nFor the MEMS vibrator 201, a MEMS vibrator utilizing change in electrostatic capacitance due to mechanical displacement is generally used in many cases. As the simplest structure, a cantilever structure can be cited as a representative example. FIG. 15 shows an exemplary appearance of the MEMS vibrator 201 using this structure; and FIGS. 16A to 16C show the configuration of the MEMS vibrator 201 in detail.\nAs shown in FIGS. 15 and 16A, the MEMS vibrator 201 is composed of a fixed electrode 101 formed on a substrate 100 and a movable electrode 102 formed apart from the fixed electrode 101 with a constant gap. The movable electrode 102 is composed of a fixed portion 104 formed on the substrate 100, a movable portion (beam) 103 provided to face the fixed electrode 101 and capable of vibrating, and a supporting portion 105 coupling the movable portion 103 with the fixed portion 104 and supporting them (refer to JP-A-2008-153817 (FIG. 2), for example).\nA method of manufacturing the MEMS vibrator 201 is as follows.\nAs shown in FIGS. 16B and 16C, a silicon oxide film 111 is formed on a Si substrate 110, and a silicon nitride film 112 is formed on the silicon oxide film 111. Next, the fixed electrode 101 formed of a first polysilicon film is formed on the silicon nitride film 112, and an insulating film (not shown) covering the fixed electrode 101 is formed. Next, the movable electrode 102 formed of a second polysilicon film is formed on the insulating film. Next, the insulating film around the movable electrode 102 is removed. In this manner, the MEMS vibrator 201 is manufactured.\nIn the MEMS vibrator 201, the movable portion 103 vibrates in a thickness direction of the substrate 100. Since the gap between the movable portion 103 and the fixed electrode 101 can be set narrow, the capacitance change rate relative to the fluctuation of the movable portion 103 is high, so that the MEMS vibrator 201 can obtain high sensitive resonance characteristics.\nThe MEMS vibrator 201 is manufactured by a so-called surface MEMS process in which the fixed electrode 101 and the movable electrode 102 are formed by depositing polysilicon or the like by a CVD (chemical vapor deposition) method. Since the MEMS vibrator 201 can be formed by the relatively simple process, it has such an advantage that the MEMS vibrator can be manufactured at low cost.\nOn the other hand, the MEMS vibrator 201 has the following disadvantage.\nAs described above, since the fixed electrode 101 and the movable electrode 102 are formed by depositing polysilicon or the like by a CVD method, the film formation by deposition results in great variation in film thickness, and the tolerance usually reaches up to ±10% in some cases. In a vibration mode, since the film thickness of the fixed electrode 101 and the movable electrode 102 serves as a factor in determining the resonant frequency of the vibrator, the variation in film thickness is directly linked to variation in resonant frequency. Therefore, in the fixed electrode 101 and the movable electrode 102 formed by depositing polysilicon or the like, there is a problem that it is difficult to obtain relatively high frequency accuracy compared to a vibrator which vibrates in a horizontal direction of a substrate.\nMeans for securing resonant frequency accuracy by managing and selecting the film thickness, optimizing the beam strength, or the like is also conceivable. In view of mass production, however, the process control has its limit, and an increase in manufacturing cost due to the addition of processes therefor is of concern.\nFIG. 17 shows the relation between variation in resonant frequency caused by variation in the fixed electrode and movable electrode of a MEMS vibrator and the target accuracy (tolerance) of resonant frequency required for the MEMS vibrator as a product. However, the degree of variation in film thickness varies depending on the performance of a film-forming apparatus, film-forming conditions, and the like, and therefore, the variation in resonant frequency shown in FIG. 17 is illustrative only. The target accuracy also varies depending on products, and therefore, the target accuracy shown in FIG. 17 is illustrative only.\nFor example, when the specification value for the resonant frequency and frequency accuracy of a MEMS vibrator is 20 MHz±100 kHz (±0.5%), the resonant frequency of the vibrator has to fall within a range from 19.9 MHz to 20.1 MHz. For realizing this, the variation in the film thickness of the MEMS vibrator has to be a desired value or less (for example, within ±0.05 μm).\nAs described above, however, if the variation in the film thickness of the MEMS vibrator results in, for example, four times the desired value (for example, ±0.2 μm) due to the performance of a film-forming apparatus, or film-forming conditions, it is presumed that the resonant frequency of MEMS vibrators which are mass-produced under this condition varies from 19.6 MHz to 20.4 MHz (±400 kHz=±2%).\nAs described above, when the variation in resonant frequency caused by the variation in film thickness is great compared to the target accuracy required as a product, the related-art manufacturing method cannot obtain a MEMS vibrator having a resonant frequency within the target accuracy with high yield."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of protein engineering. More particularly, it concerns antibodies that immunologically bind with improved affinity to a Bacillus anthracis antigen.\n2. Description of Related Art\nThe isolation of polypeptides that either bind to ligands with high affinity and specificity or catalyze the enzymatic conversion of a reactant (substrate) into a desired product is a key process in biotechnology. Ligand-binding polypeptides, including proteins and enzymes with a desired substrate specificity can be isolated from large libraries of mutants, provided that a suitable screening method is available. Small protein libraries composed of 103-105 distinct mutants can be screened by first growing each clone separately and then using a conventional assay for detecting clones that exhibit specific binding. For example, individual clones expressing different protein mutants can be grown in microtiter well plates or separate colonies on semisolid media such as agar plates. To detect binding the cells are lysed to release the proteins and the lysates are transferred to nylon filters, which are then probed using radiolabeled or fluorescently labeled ligands (DeWildt et al. 2000). However, even with robotic automation and digital image systems for detecting binding in high density arrays, it is not feasible to screen large libraries consisting of tens of millions or billions of clones. The screening of libraries of that size is required for the de novo isolation of enzymes or protein binders that have affinities in the subnanomolar range.\nThe screening of very large protein libraries has been accomplished by a variety of techniques that rely on the display of proteins on the surface of viruses or cells. The underlying premise of display technologies is that proteins engineered to be anchored on the external surface of biological particles (i.e., cells or viruses) are directly accessible for binding to ligands without the need for lysing the cells. Viruses or cells displaying proteins with affinity for a ligand can be isolated in a variety of ways including sequential adsorption/desorption form immobilized ligand, by magnetic separations or by flow cytometry (U.S. Pat. Nos. 5,223,409, 5,837,500, Georgiou et al. 1997, Shusta et al. 1999).\nThe most widely used display technology for protein library screening applications is phage display. Phage display is a well-established and powerful technique for the discovery of proteins that bind to specific ligands and for the engineering of binding affinity and specificity (Rodi and Makowski, 1999). In phage display, a gene of interest is fused in-frame to phage genes encoding surface-exposed proteins, most commonly pIII. The gene fusions are translated into chimeric proteins in which the two domains fold independently. Phage displaying a protein with binding affinity for a ligand can be readily enriched by selective adsorption onto immobilized ligand, a process known as “panning”. The bound phage is desorbed from the surface, usually by acid elution, and amplified through infection of E. coli cells. Usually, 4-6 rounds of panning and amplification are sufficient to select for phage displaying specific polypeptides, even from very large libraries with diversities up to 1010. Several variations of phage display for the rapid enrichment of clones displaying tightly binding polypeptides have been developed (Duenas and Borrebaeck, 1994; Malmborg et al., 1996; Kjaer et al., 1998; Burioni et al., 1998; Levitan, 1998; Mutuberria et al., 1999; Johns et al., 2000).\nOne of the most significant applications of phage display technology has been the isolation of high affinity antibodies (Dall'Acqua and Carter, 1998; Hudson et al., 1998; Hoogenboom et al., 1998; Maynard and Georgiou, 2000). Very large and structurally diverse libraries of scFv or FAB fragments have been constructed and have been used successfully for the in vitro isolation of antibodies to a multitude of both synthetic and natural antigens (Griffith et al., 1994; Vaughan et al., 1996; Sheets et al., 1998; Pini et al., 1998; de Haard et al., 1999; Knappik et al., 2000; Sblattero and Bradbury, 2000). Antibody fragments with improved affinity or specificity can be isolated from libraries in which a chosen antibody had been subjected to mutagenesis of either the CDRs or of the entire gene CDRs (Hawkins et al., 1992; Thompson et al., 1996; Chowdhury and Pastan, 1999). Finally, the expression characteristics of scFv, notorious for their poor solubility, have also been improved by phage display of mutant libraries (Deng et al., 1994; Coia et al., 1997).\nHowever, several spectacular successes notwithstanding, the screening of phage-displayed libraries can be complicated by a number of factors. First, phage display imposes minimal selection for proper expression in bacteria by virtue of the low expression levels of antibody fragment gene III fusion necessary to allow phage assembly and yet sustain cell growth (Krebber et al., 1996, 1997). As a result, the clones isolated after several rounds of panning are frequently difficult to produce on a preparative scale in E. coli. Second, although phage displayed proteins may bind a ligand, in some cases their un-fused soluble counterparts may not (Griep et al., 1999). Third, the isolation of ligand-binding proteins and more specifically antibodies having high binding affinities can be complicated by avidity effects by virtue of the need for gene III protein to be present at around 5 copies per virion to complete phage assembly. Even with systems that result in predominantly monovalent protein display, there is nearly always a small fraction of clones that contain multiple copies of the protein. Such clones bind to the immobilized surface more tightly and are enriched relative to monovalent phage with higher affinities (Deng et al., 1995; MacKenzie et al., 1996, 1998). Fourth, theoretical analysis aside (Levitan, 1998), panning is still a “black box” process in that the effects of experimental conditions, for example the stringency of washing steps to remove weakly or non-specifically bound phage, can only be determined by trial and error based on the final outcome of the experiment. Finally, even though pIII and to a lesser extent the other proteins of the phage coat are generally tolerant to the fusion of heterologous polypeptides, the need to be incorporated into the phage biogenesis process imposes biological constraints that can limit library diversity. Therefore, there is a great need in the art for techniques capable of overcoming these limitations.\nProtein libraries have also been displayed on the surface of bacteria, fungi, or higher cells. Cell displayed libraries are typically screened by flow cytometry (Georgiou et al. 1997, Daugherty et al. 2000). However, just as in phage display, the protein has to be engineered for expression on the outer cell surface. This imposes several potential limitations. For example, the requirement for display of the protein on the surface of a cell imposes biological constraints that limit the diversity of the proteins and protein mutants that can be screened. Also, complex proteins consisting of several polypeptide chains cannot be readily displayed on the surface of bacteria, filamentous phages or yeast. As such, there is a great need in the art for technology which circumvents all the above limitations and provides an entirety novel means for the screening of very large polypeptide libraries.\nThe use of techniques allowing creation of antibodies with improved affinities to anthrax antigens in particular is needed. Anthrax (B. anthracis) is a zoonotic soil organism endemic to many parts of the world. Infection through the inhalation of the heat resistant spores of the Gram positive bacterium, B. anthracis, results in up to 80% mortality rate if left untreated (Shafazand, 1999). In fact, B. anthracis was one of the first biological warfare agents to be developed and continues to be perceived as a major threat. While vaccine strains have been developed, widespread use is neither available nor recommended by the CDC.\nFollowing inhalation, the B. anthracis spores germinate in the alveolar macrophages and migrate to lymph nodes where they multiply and enter the bloodstream, quickly reaching 107-108 organisms per milliliter of blood (Dixon et al., 1999). The vegetative bacteria excrete a tripartite exotoxin that is responsible for the etiology of the disease. The toxin is an 83 kDa polypeptide, protective antigen (PA), that binds to a recently identified receptor on the surface of macrophages (Bradley et al., 2001). Identification of improved antibodies specific for anthrax antigens would offer the potential for new and improved therapeutic and diagnostic regimens and thus represent an important advance."}
{"text": "In today's day and age, the vast majority of businesses retain extensive amounts of data regarding various aspects of their operations, such as inventories, customers, products, etc. Data about entities, such as people, products, parts or anything else may be stored in digital format in a data store such as a computer database. These computer databases permit the data about an entity to be accessed rapidly and permit the data to be cross-referenced to other relevant pieces of data about the same entity. The databases also permit a person to query the database to find data records pertaining to a particular entity, such that data records from various data stores pertaining to the same entity may be associated with one another.\nA data store, however, has several limitations which may limit the ability to find the correct data about an entity within the data store. The actual data within the data store is only as accurate as the person who entered the data, or an original data source. Thus, a mistake in the entry of the data into the data store may cause a search for data about an entity in the database to miss relevant data about the entity because, for example, a last name of a person was misspelled or a social security number was entered incorrectly, etc. A whole host of these types of problems may be imagined: two separate record for an entity that already has a record within the database may be created such that several data records may contain information about the same entity, but, for example, the names or identification numbers contained in the two data records may be different so that it may be difficult to associate the data records referring to the same entity with one other.\nFor a business that operates one or more data stores containing a large number of data records, the ability to locate relevant information about a particular entity within and among the respective databases is very important, but not easily obtained. Once again, any mistake in the entry of data (including without limitation the creation of more than one data record for the same entity) at any information source may cause relevant data to be missed when the data for a particular entity is searched for in the database. In addition, in cases involving multiple information sources, each of the information sources may have slightly different data syntax or formats which may further complicate the process of finding data among the databases. An example of the need to properly identify an entity referred to in a data record and to locate all data records relating to an entity in the health care field is one in which a number of different hospitals associated with a particular health care organization may have one or more information sources containing information about their patient, and a health care organization collects the information from each of the hospitals into a master database. It is necessary to link data records from all of the information sources pertaining to the same patient to enable searching for information for a particular patient in all of the hospital records.\nThere are several problems which limit the ability to find all of the relevant data about an entity in such a database. Multiple data records may exist for a particular entity as a result of separate data records received from one or more information sources, which leads to a problem that can be called data fragmentation. In the case of data fragmentation, a query of the master database may not retrieve all of the relevant information about a particular entity. In addition, as described above, the query may miss some relevant information about an entity due to a typographical error made during data entry, which leads to the problem of data inaccessibility. In addition, a large database may contain data records which appear to be identical, such as a plurality of records for people with the last name of Smith and the first name of Jim. A query of the database will retrieve all of these data records and a person who made the query to the database may often choose, at random, one of the data records retrieved which may be the wrong data record. The person may not often typically attempt to determine which of the records is appropriate. This can lead to the data records for the wrong entity being retrieved even when the correct data records are available. These problems limit the ability to locate the information for a particular entity within the database.\nTo reduce the amount of data that must be reviewed, and prevent the user from picking the wrong data record, it is also desirable to identify and associate data records from the various information sources that may contain information about the same entity. There are conventional systems that locate duplicate data records within a database and delete those duplicate data records, but these systems may only locate data records which are substantially identical to each other. Thus, these conventional systems cannot determine if two data records, with, for example, slightly different last names, nevertheless contain information about the same entity. In addition, these conventional systems do not attempt to index data records from a plurality of different information sources, locate data records within the one or more information sources containing information about the same entity, and link those data records together. Consequently, it would be desirable to be able to associate data records from a plurality of information sources which pertain to the same entity, despite discrepancies between attributes of these data records and be able to assemble and present information from these various data records in a cohesive manner. In practice, however, it can be extremely difficult to provide an accurate, consolidated view of information from a plurality of information sources. The challenge can be even greater in some cases where data records from the plurality of information sources contain more than one language."}
{"text": "The present invention relates, generally, to artificial neural networks and, more particularly, to a method and circuit for performing a diagnostic which guarantees the integrity of the neurons forming the artificial neural network.\nArtificial neural networks (ANNs) are used with increased frequency in applications where no mathematical algorithm can describe the problem to be solved. Moreover, they have proven to be highly successful in classifying and recognizing objects. ANNs give excellent results in such cases because of their ability to learn by example and their proficiency of generalizing in order to respond to an input vector that was never introduced before. Thus far, most ANNs have been implemented in software and only a few in hardware.\nSeveral neurons and artificial neural network architectures implemented in semiconductor chips are described in the following related patents:\nU.S. Pat. No. 5,621,863 xe2x80x9cNeuron Circuitxe2x80x9d, issued on Apr. 15, 1997 to Boulet et al.;\nU.S. Pat. No. 5,701,397 xe2x80x9cCircuit for Pre-charging a Free Neuron Circuitxe2x80x9d, issued on Dec. 23, 1997 to Steimle, et al.;\nU.S. Pat. No. 5,710,869 xe2x80x9cDaisy Chain Circuit for Serial Connection of Neuron Circuitsxe2x80x9d, issued on Jan. 20, 1998 to Godefroy, et al.;\nU.S. Pat. No. 5,717,832 xe2x80x9cNeural Semiconductor Chip and Neural Networks Incorporated Thereinxe2x80x9d, issued on Feb. 18, 1998 to Boulet, et al.; and\nU.S. Pat. No. 5,740,326 xe2x80x9cCircuit for Searching/Sorting Data in Neural Networksxe2x80x9d, issued on Apr. 14, 1998 to Boulet, et al.,\nall of which are incorporated herein by reference.\nThe semiconductor chips can be advantageously be implemented using chips known under the name ZISC (ZISC is a registered Trademark of the IBM Corporation). In the aforementioned ZISC chips, the ANNs are based on mapping an input space, in accordance to a well known Region of Influence (ROI) algorithm.\nWith reference to U.S. Pat. No. 5,710,869, the base ZISC chip architecture is depicted in FIGS. 4A and 4B. As apparent from the drawings, a plurality of neuron circuits referenced 11-1 to 11-n are fed in parallel by four buses labeled DATA-BUS, CAT-BUS, MaxIF/MinIF-BUS and NO/CXT-BUS which respectively transmit data, category information, maximum/minimum values of the field of influence and xe2x80x98normal contextxe2x80x99 values. After being processed, the data is outputted to a common communication bus labeled COM*-BUS via OR gate 12. The detailed architecture of a neuron is shown in FIG. 5 of the above mentioned patent. During the recognition phase, the multi-norm distance evaluator circuit calculates the distance D between the input vector and the prototype stored therein. A xe2x80x98distance compare circuitxe2x80x99 compares D with the prototype actual influence field (AIF) to generate a first and second comparison signal. An identification circuit processes the comparison signals and the local category signal to generate signals that represent the neuron\"\"s response to the input vector. A minimum distance determination circuit establishes the minimum distance Dmin between the distances calculated by each neuron of the ANN. This circuit is also subsequently used to search and sort categories. For each neuron, a daisy chain circuit is serially connected to the daisy chain circuit of two adjacent neurons, thereby linking the neurons forming the ANN. The daisy chain circuit is essential in determining the state of the neuron, i.e., whether free or engaged. Finally, a context circuit enables (or inhibits) the neuron from participating alongside other neurons during the recognition phase by comparing the local context of each neuron to the global context.\nDuring the learning phase, a xe2x80x98write category operationxe2x80x99only engages the first free neuron of the ANN. As mentioned previously, the daisy chain circuit distinguishes between the two possible states of the neuron (engaged or free) and identifies the first free xe2x80x98or ready to learnxe2x80x99neuron of the ANN, based on the respective values of its input and output signals. In the following description, a neuron will be considered as having three potential states (first free, free and engaged). The free neuron is the only neuron having these signals in a complementary state. The daisy chain circuit is described in detail in U.S. Pat. No. 5,710,869 with reference to FIGS. 5 and 31. In order to facilitate the understanding of the present invention, this drawing has been duplicated herein as FIG. 1 but with different numerals.\nReferring to the aforementioned FIG. 1, the conventional daisy chain circuit now referenced 10 includes a 1-bit register 11 controlled by a store enable signal ST which is activated at initialization or at the end of the learning phase in which a new neuron is engaged. The 1-bit register 11 input is connected to the output of a 2-way AND gate 12 receiving signals DCI and RESET-. The output of 1-bit register 11 is connected to the input of a second 2-way OR gate 13 while the ALL signal is applied to the other input. The ALL signal enables all the daisy chain circuits of the ANN at initialization. DCO is the signal which is outputted from OR gate 13. Daisy chain circuit 10 further includes a control logic circuit 14 which consists of a 2-way XOR (exclusive OR) gate 15 and a 2-way AND gate 16. In XOR gate 15, DCI and DCO signals are compared to generate control signal RS (Ready to Store) that will be subsequently used as a gating signal for different circuits of each neuron to identify the first free or xe2x80x98ready to learnxe2x80x99neuron in the ANN. On the other hand, AND gate 16 generates a CO signal from the DCO and NS signals. The NS signal is generated by the matching circuit of the neuron (referenced 150 in FIGS. 5 and 16 of U.S. Pat. No. 5,710,869).\nAt initialization, all the 1-bit registers are forced to xe2x80x980xe2x80x99 by signal RESET-. The DCI input of the daisy chain of the first neuron in the ANN is set to xe2x80x981xe2x80x99 so that this neuron is xe2x80x98ready to learnxe2x80x99 (i.e., the first free neuron).\nDuring the learning phase, the 1-bit register 11 of this free neuron is set to xe2x80x981xe2x80x99 when the ST signal is activated, xe2x80x9cengagingxe2x80x9d this neuron (i.e., its output is at xe2x80x981xe2x80x99). Once engaged, the next neuron in the ANN becomes the first free neuron. The signal RS, which is a combination of both DCI and DCO signals, enables the loading of the input vector components into the R/W (read/write) memory circuit (referenced 250 in FIG. 5 of U.S. Pat. No. 5,710,869). It also determines the category specified by the user into the category register, the global context value into the local context register and, finally, the minimum distance Dmin determined by the determination circuit (referenced 500 in FIGS. 5 and 8 of the aforementioned patent) which is stored in the AIF register as the AIF value of the first free neuron. At the end of the learning process, the DCI signal is latched into the 1-bit register 11 for each neuron to define which neuron is the first free neuron of the ANN.\nDuring the recognition mode, only the neurons which are engaged and for which the local context matches the global context of all the engaged neurons (when the global context equals xe2x80x980xe2x80x99) participate in the recognition process. It is the role of the CO signal to select the distance evaluator circuit (referenced 200 in FIGS. 5 and 19 of U.S. Pat. No. 5,710,869) and the minimum distance determination circuit of these neurons to make them operative.\nUnfortunately, the daisy chain circuit in ZISC chips does not include an adequate on-chip testability capability which would easily allow to test the integrity of the neurons forming the ANN not only during the setup but also when an application is running, without complicating the existing software and negatively impacting the overall performance of the ANN.\nIt is therefore a primary object of the present invention to provide a method and a circuit that perform an integrity diagnostic of all the neurons forming an artificial neural network based upon mapping the input space.\nIt is another object of the present invention to provide a method and a circuit that perform an integrity diagnostic of all the neurons forming an artificial neural network based upon mapping the input space which can simultaneously test all the neurons in real time.\nIt is another object of the present invention to provide a method and a circuit that perform an integrity diagnostic of all the neurons forming an artificial neural network based upon mapping the input space at the cost of only a few logic gates.\nThe accomplishments of these and other related objects are achieved by the method of the present invention. The diagnostic method engages simultaneously all the neurons of the ANN network of the same prototype.\nIn a first aspect of the invention, the method includes the steps of xe2x80x98learningxe2x80x99all the neurons of the ANN of the same prototype (vector components, category, context and AIF). A first test vector having. e.g., the same context, is first loaded into the ANN. All the neurons compute their own distance with a first test vector (this distance should be the same if all the neurons were good). If at least one neuron fails, either the distance, the category, the context or the AIF value will be different from a corresponding expected value. The neuron giving at least one different answer can easily be isolated using the minimum distance determination circuit and flagged according to its particular context, i.e., a value identifying a fail. As long as this context is not selected by the global context, failing neurons do not contribute to the learning and recognition processes. The diagnostic method allows a fast integrity diagnostic of the neurons of an ANN based on mapping an input space. When implemented with ZISC chips, the diagnostic method also reveals itself to be inexpensive since it only requires a few additional logic gates."}
{"text": "The present invention relates to a direct viewing type liquid crystal display device which is employed in a liquid crystal television set, a liquid crystal display for a computer and the like.\nIn recent years, the technical progress of a direct viewing type liquid crystal display device is significant especially in a device employing a color display element. There are many displays having a display quality which is comparable to that of a CRT. In the black-white display, until several years before, the main stream of technology had been a reflecting type liquid crystal display element which did not employ a backlight. However, currently, almost all of them are replaced by a transmitting type liquid crystal display element in use of a backlight even in the black-white display. In the color liquid crystal display, the display can not work without a backlight, and the backlight is an indispensable device in the direct viewing type liquid crystal display device.\nIn a so-called xe2x80x9cnotebook type personal computerxe2x80x9d which has come in use in recent years, the portability is important, and therefore, driving thereof by a battery is indispensable. However, currently, the time capable of driving the device without charging the battery is only several hours, and does not reach a level wherein a day\"\"s operation can continuously be performed. From this viewpoint, the prolongation of the time of continuous use is extremely important. Especially, the illumination device is a device consuming much power in that system, and reducing the power consumption of the illumination device is of great significance.\nIn the mean time, there is a specific contrast ratio distribution in accordance with the viewing angle in a liquid crystal display element employed in the notebook type personal computer. A representative example is shown in FIG. 5 in case of a super twisted nematic liquid crystal display element. FIG. 5 shows the representative viewing cone of the super twisted nematic liquid crystal display element and a diagram designating a change in the contrast ratio of a section on a horizontal line passing through the center of the contrast ratio of the viewing cone. According to FIG. 5, the viewing angle of the liquid crystal display element is widened from a direction perpendicular to the image plane substantially by 40xc2x0 through 50xc2x0, and it is revealed that there is a region having an especially high contrast ratio in the vicinity of the center.\nIn an actual use, the device is often set such that the contrast ratio is maximized on the image plane viewed from an operator in case of the notebook type personal computer. Accordingly, the illuminating efficiency can substantially be promoted, when the maximum brightness of illumination is formed in the viewing direction (mostly a direction perpendicular to an image plane or a direction a little deviated therefrom) which causes the maximum contrast ratio.\nThe viewing angle of the liquid crystal display element employed in the notebook type personal computer is widened substantially to not less than 40xc2x0 through 50xc2x0, and investigations are being performed for widening the viewing angle. Therefore, it is important to adjust the light direction distribution of the illumination device so that the display can be viewed from a more or less oblique direction.\nThe adapting of the brightness distribution of the illumination device to this contrast ratio distribution in such a manner is significant as a means for promoting substantial brightness.\nThe color liquid crystal display device is grossly classified into two systems, namely, a twisted nematic (TN) liquid crystal display device by the active matrix driving using TFTs and a super twisted nematic (STN) liquid crystal display device by the multiplex driving. Both are provided with a construction wherein polarizing sheets are disposed on the light-incident side and the light-emitting side of an element wherein the liquid crystal layer is hold by glass substrates, and the liquid crystal display system is operated by modulating the polarization state of an incident linearly polarized light.\nHowever, the direction of polarization of the incident light of the conventional liquid crystal display element, is not uniform due to the randomly polarized light. Therefore, more than half of the incident light is absorbed by the polarizing sheet disposed on the light-incident side of the display element, and does not substantially contribute to the illuminating light. It is proposed as a structure for reusing the light to be absorbed by the polarizing sheet, wherein a polarized light separator is interposed between a light source lamp and a liquid crystal display device in a projecting type liquid crystal display element, for separating a non-polarized light to mutually orthogonal polarized lights, one polarized light is directly emitted from the polarized light separator, and the other polarized light is converged to the light source lamp and is reused as a light source light (Japanese Unexamined Patent Publication No. 184429/1992).\nHowever, this method is performed on the premise of a projector (projecting type), and it is necessary to provide a sufficient distance between the light source and the polarized light separator. Further, the device effectively functions as an illumination for the projecting type liquid crystal display element, only when the light is a considerably collimated parallel ray. Accordingly, the device is not suitable to adopt as the illumination for a direct viewing type display element wherein thinning thereof is an indispensable condition, and the brightness distribution of the illumination device should adapt to the contrast ratio distribution of the liquid crystal display.\nFurther, it is proposed that a prism array is interposed between a light source for illumination and a display element, as a means for converging light in a direction perpendicular to a display face. However, according to this means, the brightness in the direction perpendicular to the display face is promoted by narrowing the direction of an illuminating light in a specified range, and therefore, the light direction distribution of the illuminating light is narrowed. Further, the brightness in the direction perpendicular to the display face is not sufficient even by this means. Accordingly, an illuminance distribution suitable for the direct viewing type liquid crystal display element can not be provided.\nIn this invention, attention has been paid to the fact wherein only a polarized light having a specific direction of polarization contributes to the promotion of the illuminance of the directly viewing type liquid crystal display element, and a light in a specific direction is selectively converted its polarization direction, among polarized lights which do not contribute to the promotion of the illuminance of the liquid crystal display element by themselves. In this way, the intensity of light can be enhanced in the specific direction while maintaining a wide light direction distribution, with respect to a polarized light having a direction of polarization which can contribute to the promotion of the illuminance, and the light distribution is suitable for the directly viewing type liquid crystal display.\nAccording to a first aspect of the present invention, there is provided an illumination device for a direct viewing type display element comprising:\na flat light guide;\na light source set such that light is incident on a side portion of said flat light guide;\na polarized light separating sheet set on a first side of a light emitting side of the flat light guide for transmitting a p polarized light component and reflecting at least a portion of an s polarized light component with respect to a light ray substantially having a predetermined direction of incidence; and\na light reflecting sheet disposed on a second side opposite to said light emitting side of the flat light guide in parallel with the light emitting site.\nAccording to a second aspect of the present invention, there is provided the illumination device for a direct viewing type display element according to the first aspect, wherein the polarized light separating sheet is comprising a multi-layered structure wherein light transmitting media having a relatively large refractive index and light transmitting media having a relatively small refractive index are laminated.\nAccording to a third aspect of the present invention, there is provided the illumination device for a direct viewing type display element according to the first aspect, wherein the polarized light separating sheet comprises a transparent supporter and at least one dielectric thin film laminated on said transparent supporter having a thickness which is equal to or smaller than a wavelength of visible light.\nAccording to a fourth aspect of the present invention, there is provided the illumination device for a direct viewing type display element according to the first aspect, wherein the polarized light separating sheet comprises a plurality of laminated transparent polymer layers having different refractive indices.\nAccording to a fifth aspect of the present invention, there is provided a liquid crystal display device, wherein the illumination device according to Claim 1 is disposed on a rear side of a direct viewing type liquid crystal display element such that a principle polarization direction of emitted light from the illumination device substantially agrees with a direction of an optical axis of polarization of a polarizing sheet on a light-incident side of a liquid crystal display element.\nAccording to a sixth aspect of the present invention, there is provided the liquid crystal display device according to the fifth aspect further comprising:\na light reflecting means disposed between the polarized light separating sheet and the liquid crystal display element for deflecting a direction of a light ray maximizing a light intensity among light distributing directions to a direction substantially perpendicular to a display face of the liquid crystal display element.\nAccording to a seventh aspect of the present invention, there is provided the liquid crystal display device according to the fifth aspect, further comprising:\na means for rotating polarization direction disposed between the illumination device and the liquid crystal display element for rotating the principle direction of emitted light.\nAccording to an eighth aspect of the present invention, there is provided an illumination device for a direct viewing type display device comprising:\na flat light emitting means for emitting a diffused light including a first polarized light component having a first direction of polarization and a second polarized light component having a second direction of polarization perpendicular to the first direction of polarization; component and\na polarized light converting means disposed in an optical path of a light emitted from said flat light emitting means for emitting said first polarized light component and for emitting at least a portion of said second polarized light component after converting said portion of the second polarized light component into the first polarized light component selectively with respect to a light ray substantially having a light direction maximizing a brightness thereof."}
{"text": "The present invention relates to an assembly for holding and supporting at least one beverage container within a vehicle. More particularly, this invention relates to an assembly integrally constructed with a vehicle arm rest and adapted to hold and support multiple beverage containers. Generally, these assemblies are known as cup holders.\nWhen traveling in a motor vehicle, whether it is a short cross town commute or a trip of a longer duration, the occupants of the vehicle often find it enjoyable to consume a beverage while in route. For example, morning commuters often carry with them a hot beverage in a cup from the house or a convenient drive-thru service. At other times the beverage may be contained in a bottle or can. In all of these situations, while the vehicle is moving the container can be easily upset if it is placed on the dash, floor, seat or center console. For obvious reasons, it is desirable to provide a stable support or holder for beverage containers when they are used within a vehicle.\nOver the years, a number of devices have been developed to prevent a beverage container from spilling its contents within a vehicle. One such device is known as the travel mug. Travel mugs generally come in two varieties, one of which has a narrow opening and a wide base and the other includes a lid and sometimes a flanged bottom which can be slid into a base mounted to the dash or center console of the vehicle. Both of these mechanisms have limitations in that they require the use of a dedicated container, exclusive of other containers. Since you must always have the travel mug in your car or with you to reap its benefits, these beverage containers can be inconvenient.\nAnother direction in which designers have moved has been to develop \"generic\" cup holders which can receive and support numerous types, styles and sizes of beverage containers. One of the first of these cup holders was a one-piece plastic \"hanger\" which included a downwardly extending tab that, during use, was positioned so as to extend into the window opening generally between the glass of the window and the interior of the door. These cup holders were rigid structures and of limited use because they were bulky, easily broken, and sometimes interfered with turning of the steering wheel or the opening of the door in certain vehicles. Cup holders have also been seen which fold down and generally operate as a tray on which the container can be placed.\nNotwithstanding the above, original equipment manufacturers (OEMs) have been developing retractable cup holders which are mounted within the vehicle. One such cupholder includes a tray which is mounted for sliding movement into a recess in the dash board of the vehicle. The tray may include one or more fixed circular apertures into which the beverage container can be received. Another type of OEM cup holder includes a pivoting tray having one or more fixed circular apertures into which the beverage container can be received. These trays often pivot from a stored position to horizontal \"use\" position in response to the opening of a cover or door. The cover often operates as a rest for the bottom of a container received within the aperture of the tray.\nAs instrument panels have become increasingly crowded with electronic and ventilation equipment, cup holder trays with fixed circular apertures have fallen into disfavor. In an effort to reduce size, numerous assemblies have been developed where one or more retaining member forms less than a complete aperture into which the container is received. These retaining members often pivot with respect to a tray-type member which defines the remaining portion for the container receiving aperture. While being slightly more complicated in construction than fixed ring cup holders, these assemblies allow the cup holder to occupy less space when stored.\nAnother current trend is to relocate the cup holder assembly from the instrument panel to another portion of the vehicle. Because it is desirable to provide a compact cup holder which occupies a minimum amount of storage space, one location which has received considerable attention in the design of cup holders is the center console or fold down arm rest, respectively found in vehicles with bucket seats or split bench seats.\nGenerally, cup holders found in a center console or arm rest (hereinafter just arm rest) consist of one of two types. In the first of these types, the cup holder pivots about a horizontal axis from a stored position, located within the arm rest, into a use position where a beverage container can be received in it. Typically, this variety of cup holder is used where the arm rest also doubles as a center storage compartment and is provided with a hinged lid. The cup holder portion itself usually only occupies the forward end of the storage compartment.\nThe second type of arm rest cup holder utilizes a drawer which slides into a recess defined in the arm rest. In order to hold at least two containers whose combined width is greater than that of a standard sized arm rest, this type of cup holder must include some features which will allow those portions of the cup holder which actually define the container receiving apertures to occupy a reduced amount of space when stored in their non-use position. In accomplishing the above, cup holders have been proposed where each aperture which receives and holds a beverage container is defined in part by two elements, one of which is movable relative to the other.\nAccordingly, one or more arms are provided so that they will laterally pivot relative to a drawer as the drawer is withdrawn from a recess in the arm rest. Generally these pivoting arms include a curved portion which cooperates with a curved recess in the drawer so that when the arms are fully extended a beverage container receiving aperture is generally defined by the two. One drawback with this type of design is that since the pivoting arms only define a portion of the beverage container receiving apertures, only one end of each pivoting arm is supported making the entire arm susceptible to instability problems. Another drawback is that these types of cup holders lack the ability to accept a beverage container having a handle, such as a coffee mug.\nWith the prior art in mind, it is an object of the present invention to provide a cup holder assembly which is mounted to slide within a recess defined in the vehicle. In particular, it is an object to provide a cup holder which is mounted to slide within a center console or arm rest, located between the driver and passenger seats of the vehicle.\nAnother object of this invention is to provide a cup holder in which the beverage container receiving aperture is wholly defined by a single element so as to offer a significant amount of stability to support a beverage container therein.\nA further object of this invention is to provide a cup holder of the above mentioned variety wherein the container receiving aperture readily accepts beverage containers with handles, such as coffee mugs.\nYet another object of this invention is to provide a cup holder which will occupy a reduced amount of space when in either its stored or use position.\nStill another object of this invention is to provide a cup holder assembly which is adapted to readily receive beverage containers exhibiting a wide variety of diameters and to support them in a stable manner."}
{"text": "As a modulation method of a digital signal, orthogonal frequency division multiplexing (OFDM) is known. The OFDM is adopted in digital television broadcasting, a mobile phone, and a wireless LAN.\nMeanwhile, a portable terminal such as a smart phone and a tablet PC can be exemplified as an electronic apparatus receiving an OFDM-modulated signal (hereinafter, referred to as the OFDM signal). Because there is a limitation in a battery capacity of the portable terminal, low consumption power is demanded when the OFDM signal is received. In an LSI mounted to the portable terminal and having a function of receiving the OFDM signal to realize the demand, an SOC solution in which an analog circuit and a digital circuit are mounted on the same silicon is frequently adopted.\nMeanwhile, as miniaturization of a semiconductor such as the LSI advances, a digital circuit in the LSI can realize a small size and low consumption power. However, in an analog circuit in the LSI, because consumption power increases or a circuit scale increases to suppress an influence such as reduction of a power-supply voltage, a benefit of the miniaturization may not be used. For this reason, in the LSI in which the SOC solution is adopted, a small size and low consumption power in the analog circuit are demanded.\nFIG. 1 illustrates an example of a configuration of a receiving device according to the related art that is mounted as the LSI to the portable terminal and receives an OFDM signal.\nA receiving device 10 includes a variable gain amplifier 11, frequency converting units 12I and 12Q, BB filters 13I and 13Q, variable step gain BB amplifiers 14I and 14Q, AD converting units (ADC) 15I and 15Q, a local oscillating unit (LO Gen.) 16, an AGC control unit (AGC CTRL LOGIC) 17, and an OFDM demodulating unit (DMD) 18. The variable gain amplifier 11 continuously amplifies an OFDM signal to be an RF signal, according to RF gain notified from the AGC control unit 17, and outputs the OFDM signal to the frequency converting units 12I and 12Q. The frequency converting unit 12I mixes the amplified OFDM signal with a signal for frequency conversion from the local oscillating unit 16, converts the RF signal into an IF signal, and outputs the IF signal to the BB filter 13I. The frequency converting unit 12Q mixes the amplified OFDM signal with a signal for frequency conversion from the local oscillating unit 16, converts the RF signal into an IF signal, and outputs the IF signal to the BB filter 13Q.\nThe BB filter 13I extracts a base band signal of an I component of the OFDM signal converted into the IF signal and outputs the base band signal to the variable step gain BB amplifier 14I. The BB filter 13Q extracts a base band signal of a Q component of the OFDM signal converted into the IF signal and outputs the base band signal to the variable step gain BB amplifier 14Q.\nThe variable step gain BB amplifier 14I gradually amplifies the base band signal of the I component of the OFDM signal, according to COMP. gain notified from the AGC control unit 17, and outputs the base band signal to the AD converting unit 15I. The variable step gain BB amplifier 14Q gradually amplifies the base band signal of the Q component of the OFDM signal, according to COMP. gain notified from the AGC control unit 17, and outputs the base band signal to the AD converting unit 15Q.\nThe AD converting unit 15I converts the base band signal of the I component of the OFDM signal into a digital signal and outputs the digital signal to the AGC control unit 17 and the OFDM demodulating unit 18. The AD converting unit 15Q converts the base band signal of the Q component of the OFDM signal into a digital signal and outputs the digital signal to the AGC control unit 17 and the OFDM demodulating unit 18.\nThe local oscillating unit 16 supplies a signal for frequency conversion to the frequency converting unit 12I. In addition, the local oscillating unit 16 shifts a phase of the signal for the frequency conversion supplied to the frequency converting unit 12I by 90 degrees and supplies the signal for the frequency conversion to the frequency converting unit 12Q.\nThe AGC control unit 17 determines the RF gain and the COMP. gain on the basis of outputs of the AD converting unit 15I and 15Q and notifies the variable gain amplifier 11 and the variable step gain BB amplifiers 14I and 14Q of the RF gain and the COMP. gain, respectively. The OFDM demodulating unit 18 demodulates a transmitted symbol on the basis of the I component and the Q component of the digitized OFDM signal and outputs the symbol to a rear step.\nIn the receiving device 10, the OFDM signal input as the RF signal is amplified by the variable gain amplifier 11 according to the RF gain and is converted into the IF signal by the frequency converting units 12I and 12Q and the base band signals of the I component and the Q component are extracted by the BB filters 13I and 13Q. In addition, the I component and the Q component of the OFDM signal are amplified by the variable step gain BB amplifiers 14I and 14Q according to the COMP. gain, are digitized by the AD converting units 15I and 15Q, and are demodulated into the symbol by the OFDM demodulating unit 18.\nAs illustrated in FIG. 1, in the receiving device 10 mounted to the portable terminal, a discrete control type amplifier in which an analog circuit is reduced and a small size and low consumption power are realized is adopted in partial components (the variable step gain BB amplifiers 14I and 14Q). However, in the case in which the discrete control type amplifier is adopted, if a steep gain variation is generated, a characteristic of the OFDM demodulation may be notably deteriorated in general.\nTherefore, a method of additionally mounting a continuous control amplifier to suppress the steep gain variation from being generated is known (for example, refer to NPL 1). However, in this method, a circuit scale and consumption power may be increased.\nTherefore, as illustrated in FIG. 2, a method of performing digital correction on an ADC output is suggested as another method of suppressing the steep gain variation from being generated (for example, refer to NPL 2).\nFIG. 2 illustrates an example of a configuration of a receiving device according to the related art that includes a digital correcting unit to perform the digital correction on the ADC output, is mounted as an LSI to a portable terminal, and receives an OFDM signal.\nA receiving device 20 includes a variable step gain amplifier 21, frequency converting units 12I and 12Q, BB filters 13I and 13Q, variable step gain BB amplifiers 14I and 14Q, AD converting units (ADC) 15I and 15Q, digital correcting units 22I and 22Q, a local oscillating unit (LO Gen.) 16, an AGC control unit (AGC CTRL LOGIC) 23, and an OFDM demodulating unit (DMD) 18.\nIn the receiving device 20, the variable gain amplifier 11 and the AGC control unit 17 of the receiving device 10 of FIG. 1 are replaced by a variable step gain amplifier 21 and an AGC control unit 23, respectively, a digital correcting unit 22I is provided between the AD converting unit 15I and the OFDM demodulating unit 18, and a digital correcting unit 22Q is provided between the AD converting unit 15Q and the OFDM demodulating unit 18.\nThe variable step gain amplifier 21 gradually amplifies an OFDM signal to be an RF signal, according to the RF gain notified from the AGC control unit 17, and outputs the OFDM signal to the frequency converting units 12I and 12Q.\nThe digital correcting unit 22I is operated to correct a gain variation generated in an output of the AD converting unit 15I by a step response (gradual amplification) in an analog circuit, according to correction gain notified from the AGC control unit 23. The digital correcting unit 22Q is operated to correct a gain variation generated in an output of the AD converting unit 15Q by a step response in an analog circuit, according to correction gain notified from the AGC control unit 23.\nThe AGC control unit 23 determines the RF gain, the COMP. gain, and the correction gain, on the basis of outputs of the digital correcting units 22I and 22Q, and notifies the variable step gain amplifier 21, the variable step gain BB amplifiers 14I and 14Q, and the digital correcting units 22I and 22Q of the RF gain, the COMP. gain, and the correction gain.\nBecause components common to the components of the receiving device 10 of FIG. 1 are denoted with the same reference numerals, explanation thereof is omitted."}
{"text": "At present, a touch display screens are provided in more and more electronic devices, e.g., an information inquiry machine in a public service hall, or a computer or a mobile phone used by a user in daily life and work. Using a touch display screen, the user can operate the electronic device by simply touching icons on the touch display screen with a finger, thereby getting rid of the keyboard and mouse operations and making a human-computer interaction more straightforward. To better meet user requirements, the touch display screen is typically provided with one or more pressure sensors for detecting touch pressure values when the user touches the touch display screen.\nStatic electricity generated by a structure in a display area or a circuit in a non-display area causes damage to the pressure sensors when the touch display screen is being manufactured or is running properly. In addition, when a high intensity current flows through a power line or a ground potential line electrically connected to a signal input terminal of a pressure sensor or when the touch display screen is subjected to electrostatic testing, the pressure sensor is affected by the corresponding static electricity, causing damage or performance deterioration to the pressure sensor."}
{"text": "This invention pertains to a food packaging apparatus. More particularly, this invention pertains to a device for defining a loop handle in a tubular covering, such as an expandable mesh covering, for a food product, such as, but not limited to a meat product.\nIt is known in the food packaging industry to vacuum wrap consumer portions of food products, especially meat products. It is also known in the art to use expanded mesh coverings for food product during processing. To facilitate carrying, it has become customary to provide a tubular cover, such as an expanded plastic mesh cover for the vacuum wrapped food product. Further, this tubular mesh cover has a loop in one end in order to define a handle. Presently, formation of this loop handle is accomplished manually by hand. In this regard, a worker will bag the food product in the tubular expanded mesh covering, gather, or compress, an extended portion back on itself and then clip the loop, with a conventional clipping apparatus, in order to secure the loop. Conventional clipping machines attach two clips, one to secure the bottom of the next covering to be formed and one to secure the loop handle. The conventional clipping machine then cuts the tubular covering between these two clips. Manual loop formation can limit the rate of production. What is needed is an apparatus for defining a handle in a tubular covering for a food product."}
{"text": "1. Technical Field\nThe present invention relates to a method of fabricating a copper indium gallium selenide (CIGS) thin film for solar cells through co-vacuum evaporation and a CIGS thin film for solar cells fabricated by the same. More particularly, the present invention relates to a CIGS thin film for solar cells fabricated using a simplified co-vacuum evaporation process without deterioration in physical properties of crystal growth of the thin film and band-gap grading by Ga distribution while reducing process costs through simplification of process steps and adjustment of evaporation elements.\n2. Description of the Related Art\nRecently, importance on development of next-generation clean energy sources has increased due to the depletion of fossil fuel reserves. Thereamong, a solar cell is a device that converts solar energy directly into electricity. Solar cells can serve as an energy source to solve energy problems in the future, since they do not emit pollutants and rely upon the sun.\nA solar cell may be divided into a variety of kinds depending upon materials used in a light-absorption layer, and the most currently available solar cell is a silicon solar cell. However, since silicon prices have been rising due to shortage of high purity silicon, a thin film type solar cell is drawing attention. The thin film type solar cell is fabricated to a thin thickness and thus contributes to less material consumption and light weight, thereby providing a wide application range. Studies have been actively made as to amorphous silicon, CdTe, CIS or CIGS as a material for such a thin film solar cell.\nCIS or CIGS is a I-III-VI compound semiconductor material and has higher conversion efficiency than any other thin-film solar cell materials experimentally prepared. Particularly, since the CIS or CIGS can be fabricated to a thin thickness of 10 microns or less and can be stably operated even after long term use, the CIS or CIGS is expected to be a low-cost, high-efficiency solar cell material as an alternative to silicon.\nParticularly, a CIS thin film is a direct transition type semiconductor that can be made thinner, is suitable for optical-conversion due to a band-gap of 1.04 eV, and exhibits high optical-absorption coefficient. CIGS is an alternative to CIS in which a portion of In is substituted with Ga or Se is substituted with S in order to improve low open circuit voltage of CIS.\nA CIGS thin film can be generally fabricated by vacuum deposition or non-vacuum coating. Particularly, vacuum deposition may include co-evaporation, in-line evaporation, a two-step process (precursor-reaction), and the like. Thereamong, co-evaporation has been traditionally used to fabricate a high-efficiency CIGS thin-film solar cell. However, co-evaporation has a problem of difficulty in commercialization due to a complicated process and difficulty in fabrication of a large area solar cell. To solve this problem, a two-step (deposition/selenization) process facilitating mass production has been developed.\nAmong current CIGS deposition processes, a high efficiency technique having an efficiency of 20% or more is a three-step co-vacuum evaporation process developed by NREL wherein, in a first step, In, Ga, and Se are deposited, in a second step, Cu and Se are deposited excessively relative to a normal amount to promote crystal growth, and in a third step, In, Ga, and Se are further supplied to a thin film to adjust a composition of the CIGS thin-film. Particularly, a Cu—Se compound formed in the second step has a low melting point of about 523° C. or less, thereby causing crystal growth of the CIGS thin-film through induction of liquid-phase sintering of the compound. Thus, it has been known that a V-type distribution of Ga in the composition of the thin-film is obtained through such a three-step process, thereby providing efficiency improvement through band-gap grading.\nHowever, considering mass production, since regulation of flux of elements in each step is complex and thin-film deposition requires a long time, the three-step process has a drawback of high process cost. Therefore, there is a need for a method of fabricating a CIGS thin film for solar cells through a simplified process capable of realizing crystal growth and band-gap grading by Ga composition distribution, while simplifying process steps and significantly reducing a film-deposition time."}
{"text": "Field\nEmbodiments of the present disclosure generally relate to semiconductor processing equipment. More particularly, embodiments of the present disclosure relate to a plasma source, and an abatement system for abating compounds produced in semiconductor processes.\nDescription of the Related Art\nThe process gases used by semiconductor processing facilities include many compounds which must be abated or treated before disposal, due to regulatory requirements and environmental and safety concerns. Typically, an abatement system may be coupled to a processing chamber to abate the compounds exiting the processing chamber. The abatement system typically includes at least a plasma source. Halogen-containing plasma and gases are frequently used in etch or cleaning processes, and components of the processing chamber and the abatement system are susceptible to corrosion from the halogen-containing plasma and gases. The corrosion reduces the service life of the processing chamber components and the abatement system, and, additionally, introduces undesirable defects and contamination into the processing environment.\nAccordingly, what is needed in the art is an improved plasma source, and an abatement system for abating compounds produced in semiconductor processes."}
{"text": "This invention is directed to a cold worked stainless steel bezel for a portable electronic device.\nBy their very nature, portable electronic devices (e.g., MP3 players, cellular telephones) are carried around and subject to impacts and inadvertent blows to which static electronic devices (e.g., desk-top computers, televisions) are not subject. To protect the electronic systems of these portable devices, manufacturers have constructed impact resistant cases.\nExisting cases, however, are not always easily manufactured, aesthetically pleasing, or sufficiently resistant to impacts. Accordingly, there is a need for a hard, easily manufactured and aesthetically pleasing case for portable electronic devices."}
{"text": "Sensor assemblies can be used to measure forces or torques acting on a mechanical instrument as a result of its interaction with an object. Many sensor assemblies include strain sensors configured to measure strains experienced by a sensor body of the sensor assembly as a result of such forces or torques. The measured strains can be used to resolve one or more components of such forces or torques.\nHowever, many existing sensor assemblies are incapable of distinguishing between thermal strains resulting from temperature variations and mechanical strains resulting from forces or torques acting on the instrument. In such sensor assemblies, thermal strains may introduce significant error when measurements form strain sensors are used to determine components of forces or torques acting on the instrument.\nIn recent years some sensor assemblies have been developed that are capable of compensating for thermal strains resulting from steady state temperature variations. However, these sensor assemblies are often highly complex and require many components, resulting in a high cost of production. Furthermore, some of these sensors assemblies are limited to resolving only a subset of axial and lateral forces and axial torque independent of steady state temperature variations."}
{"text": "1. Field of the Invention\nThis invention generally relates to sealing structures. More specifically, the present invention relates fishing-reel sealing structures that seal between a spindle-form first member and a second member disposed about the outer periphery of the first member in fishing reels.\n2. Background Information\nBecause fishing reels such as spinning reels and dual-bearing reels are used in environments where they are constantly splashed with water, sealing structures are employed in fishing-reel applications for preventing water from encroaching interiorly. In particular, rings and oil-seals have been in widespread use conventionally for sealing in between relatively rotatable shaft members and their encompassing members. O-rings and oil-seals are used, for example, on the outer side of handle-carrying bearings in spinning reels or in between cover members and the spool shafts in dual-bearing reels.\nIn conventional O-ring based sealing structures, the O-rings are slightly squashed in order to form a seal. This means that the contacting surface area of the O-rings is extensive, and in particular, where the two members are relatively rotatable, the O-rings make the rotational resistance large, which tends to impair the rotational performance.\nOil seals are generally composed of a lip made of rubber, and a shape-retaining backing member made of metal. Accordingly, with oil-seal based sealing structures, because the contacting part is a peaked lip shape, it is comparatively easy to curb degradation in rotational performance. However, oil seals cannot be extended/contracted due to the backing member made of metal, and therefore fastening members and a space for fixing the oil-seal are necessary, which increases the axial space required to accommodate the oil seals.\nIn view of the above, there exists a need for sealing structure which overcomes the above mentioned problems in the prior art. This invention addresses this need in the prior art as well as other needs, which will become apparent to those skilled in the art from this disclosure.\nAn object of the present invention is to provide a structure for sealing a spindle-form first member and a second member disposed about the outer periphery thereof, to seal without having to use components or space for fastening, while curbing rotational-performance impairment.\nA fishing-reel sealing structure having is furnished with a fishing-reel, spindle-form first member, a second member, and a seal member. The second member has an inner peripheral surface disposed about the first member outer-peripherally, and is rotatable relative to the first member. The seal member, a component made of an elastic substance, seals the clearance between the two members, and has: an elastic-manufacture cylindrical portion mountable on the first member; and a sealing portion formed integrally on the outer circumferential surface of the cylindrical portion, and tapering toward the inner peripheral surface of, and whose distal edge is for contacting on, the second member.\nIn this sealing structure, the cylindrical portion of the seal member made of an elastic substance is fitted to the spindle-form first member, and bringing the tapered sealing portion into contact with the inner circumferential surface of the second member seals the clearance between the two members. Herein, the fact that the tapered sealing portion is brought into contact with the second member to seal the clearance between the two members means that the surface area of the places in contact is smaller, making degradation in rotational performance unlikely. Likewise, that fact that, like O-rings, the seal member is entirely made of an elastic substance makes it extendable/contractible, which means that, for example, it may be fastened to the first member simply by fitting its cylindrical portion into a groove. Components and space for fastening the seal member are therefore made unnecessary.\nPreferably, an annular groove is formed in a first-member outer circumferential surface that opposes the second member, and the cylindrical portion is fitted into the annular groove. In this case, the seal member may be fastened to the first member just by stretching the cylindrical portion and fitting it into the first member, and set that way shifting axially to mount the seal member in the annular groove.\nPreferably, the part of the sealing portion for contacting the second member is a lip inclined heading outward of the first member. In this case, the fact that the sealing portion is a lip inclined heading outward of the first member means that when mounting the first member from without, the second member will not reverse even though it contacts the sealing portion. Moreover, since the fore end of the sealing portion is oriented outward, though liquids enter from without, the liquids are not likely to encroach interiorly.\nPreferably, the first member is a constituent of a lever-drag mechanism in a lever-drag-type dual-bearing reel having a line-winding spool, and therein is a stationary component immovable in the axial direction of the spool; and the second member is a constituent of the lever-drag mechanism that is rotatable relative to the first member, and therein is a shifting member for shifting in the axial direction of the spool. In this case, utilizing the spool member also in a lever-drag mechanism curbs degradation in rotational performance when the spool spins freely.\nPreferably, the first member is a drag disk capable of braking the spool, and rotates linked to a spool shaft passing through the spool center; and the second member is a cover member for covering the drag disk, and rotates linked to the spool and meanwhile shifts in the axial direction together with the spool and the spool shaft. In this case, fluctuations in drag force may be controlled by sealing the clearance between the cover member and the drag disk.\nPreferably, the drag disk has a boss portion passing centrally through the cover member; the seal member is fitted to the boss portion; and a contact face for the seal to contact and a relief face the seal does not contact are formed, ranged in the spool axial direction, superficially on the inner rim of the cover member. In this case, relative axial movement between the drag disk and the cover member brings the sealing portion of the seal member into contact with/parts it from the cover member. This allows for specific conditions whereby the sealing portion does not contact the cover member so as not to impair the rotational performance.\nPreferably, the sealing portion of the sealing member: parts away from the cover member to oppose the relief face, when the drag disk in its brake-release state parts away from the spool; and comes into opposition with, to contact on, the contact face when the drag disk in its braking state contacts the spool. In this case, the fact that the during the brake-release state the sealing portion is opposite the relief face, where it is parted off the cover member, means it does not impair rotational performance when the spool spins freely, and lets fishing line be reeled out agilely though a seal member is installed. Likewise, during the braking state, the sealing portion seals positively by coming into contact with the contact face of the cover member.\nPreferably, the spool has a bobbin-trunk portion, and flange portions arranged on either end of the bobbin-trunk portion; and the outer diameter of the sealing portion is 8% or more and 50% or less of the outer diameter the flange portions. In this case, the fact that the outer diameter of the sealing portion is 50% or less that the outer diameter of the flange portions means that the contacting surface area of the sealing-member sealing portion is sufficiently small with respect to the spool, which curbs degradation in rotational performance though the seal member is brought into contact with the second member. Likewise, the fact that the outer diameter of the sealing portion is 8% or more lets strength of the spool shaft be sustained with the seal member fitted into a first member mounted encompassing the spool shaft.\nPreferably, the first member is a handle spindle fitted detachably and reattachably to a spinning-reel master-gear shaft; and the second member is a reel body rotatively carrying the master-gear shaft. In this case, impairment in rotational performance of the spinning-reel handle may be controlled; meanwhile in either case where the handle is fitted on through the left or the right, sealing either end of the master gear shaft may be accomplished with a single seal member.\nA fishing-reel sealing structure is furnished with a fishing-reel, spindle-form first member, a second member, and a seal member. The second member has an inner peripheral surface disposed about the first member outer-peripherally, and is rotatable relative to the first member. The seal member, a component made of an elastic substance, seals the clearance between the two members, and has: an elastic-manufacture cylindrical portion mountable on the second member; and a sealing portion formed integrally on the inner circumferential surface of the cylindrical portion and tapering toward the outer circumferential surface of the first member, and whose distal edge is for contacting on the first member.\nIn this sealing structure, the cylindrical portion of the seal member made of an elastic substance is fitted to the second member, and bringing the tapered sealing portion into contact with the outer circumferential surface of the first member seals the clearance between the two members. Herein, the fact that the tapered sealing portion is brought into contact with the first member to seal the clearance between the two members means that the surface area of the contacting places is smaller, making degradation in rotational performance unlikely. Likewise, that fact that, like O-rings, the seal member is entirely made of an elastic substance makes it extendable/contractible, which means that, for example, it may be fastened to the second member simply by fitting its cylindrical portion into a groove. Components and space for fastening the seal member are therefore made unnecessary.\nThese and other objects, features, aspects and advantages of the present invention will become apparent to those skilled in the art from the following detailed description, which, taken in conjunction with the annexed drawings, discloses a preferred embodiment of the present invention."}
{"text": "1. Field of the Invention:\nThe present invention relates to a lifting apparatus capable of conveying persons and/or goods from or to different levels by raising a platform from a vehicle and, more particularly, to a lifting apparatus capable of moving horizontally and raising a platform vertically.\n2. Description of the Prior Art:\nLifting apparatus capable of conveying persons and goods from or to different levels by raising or lowering a platform are used widely for assembling work, painting work or repairing work performed at high elevations on expressways and building construction sites. The typical conventional lifting apparatus employs a telescopic parallel mechanism, namely, a so-called scissors type linkage, comprising a plurality of pairs of arms with each pair of arms being pivotally joined at the middle. In order to increase the maximum lift of such a lifting apparatus, it is necessary to increase the length of the arms or to increase the number of the pairs of arms. Accordingly, such a lifting apparatus having a large lift needs a parallel mechanism consisting of many links. Therefore, the lifting platform of the lifting apparatus is located at a high level even when the parallel mechanism is collapsed, and hence it is difficult for persons to get on and off the platform and it is troublesome to load and unload the platform.\nLifting apparatus having a single extendable arm comprising a plurality of telescopically combined booms have been proposed in Japanese Patent Application Nos. 56-134487 and 56-191065.\nAll of the above-mentioned newly proposed lifting apparatus, however, inevitably need an increased number of booms, and hence so many components are necessary that troublesome manufacturing and assembling work is required and the lifting apparatus are expensive. Furthermore, the above-mentioned lifting apparatus have a large number of sliding parts for assembling the booms and arms. Since these sliding parts are provided with sliding members, such as MC nylon members, a large number of parts need to be replaced periodically requiring high inspection and maintenance costs and troublesome work.\nAnother lifting apparatus has been disclosed in Japanese Patent Application No. 59-95797. This lifting apparatus has a single extendable boom mounted on a vehicle, and the boom is telescopically moved in the manner of an expanding Z-shape as viewed from the side portion thereof. However, in this lifting apparatus it is difficult to support the platform horizontally when in a fixed elevated position since the single extendable boom is made up of a plurality of booms each of which is moved slidably in a telescopic manner and the surfaces of the booms are not even along the longitudinal axes thereof because of differences in cross section thereof. The differences in cross section between each boom thereby causes the expandable boom assembly to be unstable. Hence, it is difficult to mount a horizontal leveling mechanism on a telescopic boom to maintain the platform horizontal with respect to the chassis."}
{"text": "The present invention relates carriers for memory devices, and in particular relates to a device that both carries and protects memory cards.\nCommon to personal computers and consumer electronic components are portable memory devices, such as the well known three and one-half inch floppy diskette. That diskette includes a magnetic memory media housed within a rigid plastic case including a spring biased, metal slide gate on a top or insertion edge of the diskette case. Whenever such a diskette is inserted within a computer or consumer electronic component, such as an electronic keyboard, the slide gate slides from a closed to an open position to provide access by the computer or component to the magnetic media within the plastic case. It is well known that the three and one-half inch diskettes are commonly used for transferring data from one computer to another, for backing up a computer in order to preserve data in the event of failure of the computer, for transferring data from a computer to a consumer electronic component, and for various other common tasks. Consequently, there are many carriers available for storing the three and one-half inch diskettes.\nMost carriers for such diskettes have common features. Typically they provide for ease of access to identifying labels on the diskettes, so a user may readily select a specific diskette without searching through many diskettes to find the specific diskette. For example, a recent diskette carrier is shown in U.S. Pat. No. 5,772,020 issued on Jun. 30, 1998 to Hara. The Hara carrier includes a planar sheet having six holding cavities for six diskettes wherein a user may readily view front or label surfaces of the diskettes. Opposed, biased \"pressing strips\" extend from an edge of each cavity to apply pressure to the diskette to secure it within the cavity. The pressing strips are positioned to avoid substantial direct contact with the metal slide gate because the gate is well known to be delicate, and any excessive pressure could easily damage the slide gate, or the magnetic media under the gate within the plastic case of the diskette. U.S. Pat. No. 5,529,183 issued on Jun. 25, 1996 to Nishikawa and U.S. Pat. No. 4,676,374 issued on Jun. 30, 1987 to Wilkins are also typical of three and one-half inch diskette carriers in having resilient tabs to secure the diskette within carrier cavities, wherein the tabs are positioned adjacent opposed side edges of the data transfer port or slide gate. Consequently, the peripheral regions of cavities or diskette sleeves of the known information carriers for diskettes invariably have an exposed region, wherein the slide gate or information transfer port is not protected by the carrier against contamination or impact damage.\nRecently, a new form of information or memory storage and transfer device has been introduced into the consumer electronic field in particular. This new memory transfer device is generally referred to as a \"flash memory card\", a \"memory card\", an \"SD memory card\" or a \"MultiMedia card\". These cards are characterized as non-volatile, solid state memory devices for digital storage. For convenience, the cards will hereinafter be referred to as \"memory cards\". Known memory cards have a capacity of up to 64 megabytes, and it is anticipated that 256 megabyte memory cards will be available in the near future.\nThe memory cards are generally rectangular, being about 32 millimeters (\"mm\") long, by about 24 mm wide. Each memory card includes a front or label surface, and a back or data transfer surface. Each memory card also includes a top or insertion end which is the end of the card to be inserted into a card receiving electronic device, such as the now available \"MP3 Player\" for audio reproduction sold by PONTIS, Inc. of Schwarzenfeld, Germany.\nAdjacent the insertion end of the memory card on the back side of the card are a plurality of data transfer ports that may be isolated metal contacts. To make sure that a user does not insert the memory card into the electronic device upside down, or backwards, one corner of the insertion end of the memory card is a non-right angle, or asymmetrical corner. If a memory card was aligned so that the data transfer or back surface was hidden from a user, and the insertion edge was the top edge, then a top right corner would be the asymmetrical corner. The top left corner is symmetrical to the bottom left corner, and the two bottom corners are symmetrical to each other. By industry standards, however, the top right corner of known memory cards is asymmetrical, which may mean that the corner is formed of two angles totaling ninety degrees. For purposes herein however, the phrase \"asymmetrical corner\" is meant to include any corner in a memory storage and/or transfer device that is distinct from the remaining corners of the memory device for the purposes of correct alignment of the memory device within an electronic apparatus that receives the memory device.\nKnown carriers for rigid case memory storage devices, however, are limited to carriers for three and one-half inch types of diskettes, and those carriers do not provide adequate protection for exposed data transfer ports. If the data transfer ports on new memory cards are contaminated by liquids, salt deposits or abrasive particles, retrieval of information stored on the cards may be jeopardized, and the memory cards may actually damage the electronic devices into which the cards are inserted. Additionally, static discharge onto the data transfer ports of memory cards may destroy data stored within the cards. Accordingly there is a need for a memory card carrier device that both facilitates ease of access to and carrying of the cards and that also protects data transfer ports of the cards while the cards are stored within the carrier."}
{"text": "Unless otherwise indicated herein, the materials described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nElectric motor actuators for robotic and automation systems often require a transmission (speed reducer) in order to operate within the speed-torque requirements of a specific application and of the motor. Commonly used systems include multi-stage gearboxes, timing belts, cables with capstans, harmonic drives, and cycloid gearboxes. These systems may have shortcomings, such as too inefficient, susceptible to overload damage, heavy, and require expensive precision manufacturing. Such systems are also often too expensive for consumer products when high performance is required.\nAs an example, harmonic drive systems can be used in high performance applications where low backlash and gear-ratios greater than 50:1 are required. The harmonic drive is proprietary, heavy, inefficient, susceptible to damage from shock load, and expensive for consumer application. As another example, cable drive systems can be lightweight and efficient; however, achieving useful transmission ratios may lead to complex multi-stage designs that require high preload forces, large cable bend radii, and challenging cable management. Often it is desired to integrate a torque sensor such as a strain gauge load cell at or near the output of the transmission in order to achieve closed loop torque control. Practically, such as a case on a rotating output hub, integration of this sensor can prove challenging as the sensor wires typically rotate with the transmission output, and therefore, require cable management, adding complexity and increasing the number of failure points."}
{"text": "1. Technical Field\nThe present disclosure relates to a printed circuit board (PCB) assembly, especially to a PCB assembly for automatically detecting a computer system the PCB assembly is coupled to.\n2. Description of Related Art\nA PCB, such as a motherboard usually has multiple functions for coupling to different computer systems. The PCB may have different soft configurations to couple to these different computer systems. For example, a server motherboard may alternatively used in a 1 U (unit) server chassis or in a 2 U server chassis. A server chassis the server motherboard used is often identified by human. The corresponding configuration files may be set by adjusting a jumper on the server motherboard manually. It is time consuming to manually identify and adjust the computer system the PCB coupled in."}
{"text": "The present invention relates generally to systems for recording the picture image content of a cinematographic film in a form of a video signal on a rotating recording medium and for subsequently reproducing therefrom the video signal. More particularly, the invention relates to a system which records the picture content of a cinematographic film traveling at a speed of 24 film frames or scenes per second in a form of a television video signal of 60 fields (namely, 30 frames) per second on a rotating recording medium rotating at a speed of 900 revolutions per second in a manner such that still picture reproduction can also be carried out, and which subsequently reproduces the video signal thus recorded from the recording medium.\nIn general, in a system for recording a television video single on a rotating recording medium, the fewer the field of the video signal to be recorded per revolution of the rotating recording medium, the higher is the rotational speed at which the rotating recording medium must be rotated. In this case, the system is subjected, mechanically, to unreasonably and excessively severe conditions, whereby vibration occurs, or automatic control becomes difficult. Moreover, the recording capacity of the recording medium is reduced.\nOn the other hand, the greater the number of fields of the video signal to be recorded per revolution of the rotating recording medium, the lower is the required rotational speed of the rotating recording medium. For this reason, the system is not subjected mechanically to excessively severe conditions, and furthermore, the recording capacity of the recording medium becomes large. On the other hand, however, since the wavelength of recording on the recording medium becomes short, in the case where the signal is recorded on the recording medium as a variation of geometrical shape and is thereafter reproduced as capacitance variation or optical variation, the reproducing element becomes unable to reproduce satisfactorily.\nIn this connection, a novel \"Information signal recording and reproducing system\" assigned to the same assignee of the present application was proposed by U.S. patent application Ser. No. 785,095, filed Apr. 6, 1977, now U.S. Pat. No. 4,331,976, and is now being reduced to practice. In accordance with this system, a video signal is recorded along a track of spiral form as a geometrical shape variation, without a reproducing stylus guide groove, on a rotating recording medium having a flat surface. Since there is no guide groove in the rotating recording medium, the reproducing stylus can move smoothly from one track turn to another at the time of still-picture, slow-motion, or quick-motion reproduction. Accordingly, in spite of the use of a reproducing stylus in this system, special reproducing modes such as the still-picture reproduction can be smoothly carried out.\nHowever, in the case where still-picture reproduction is carried out by repeatedly reproducing the same track turn, if an information content signal of a larger number of fields is recorded along the same track turn, the reproduced still picture will not be one which is perfectly still but will appear unnatural and unsightly or indistinct. More specifically, particularly in the case where the information content is one of rapid motion, the position of each moving object at the starting point of the same track turn differs from that at the terminal point of the same track turn in the reproduced picture. For this reason, when the same track turn is repeatedly reproduced, each moving object repeatedly undergoes reciprocation been these two different positions, and the reproduced image appears to flicker or vibrate. Also at the time of slow-motion reproduction when the same track turn is reproduced a plurality of times and then the reproducing element (stylus) shifts to the succeeding track turn, the reproduced picture will move unnaturally.\nIn view of the various above described circumstances, we have decided that it is most desirable to record a video signal at a rate of four fields, that is, two frames, per revolution of the rotating recording medium.\nThe procedure of converting the picture image content of a cinematographic film of 24 film frames per second into a video signal of 60 fields (30 frames) per second and using this video signal for either television broadcasting or recording on a recording medium such as a magnetic tape has been practiced in the prior art. For this conversion of the picture of 24 film frames per second into a video signal of 60 fields per second, in general, a so-called 3-2 pull-down method is being used. This 3-2 pull-down method is carried out by controlling the cinematographic film feed and the shutter operation of the projector so that the first film frame is transcribed for 3 fields, the second film frame for 2 fields, the third film frame for 3 fields, the fourth film frame for 2 fields, the succeeding film frames similarly being transcribed alternately for 3 fields and 2 fields. The light image thus projected by this projector is picked up by a camera tube and thereby converted into a television video signal.\nHowever, when the video signal converted by this 3-2 pull-down method is recorded directly as it is on a rotating recording medium at the above described rate of 4 fields per revolution (corresponding to one track turn), the information of the first film frame of 3 fields and the information of the second film frame of 1 field are recorded along the first track turn; the information of the remaining fields of the second film frame and the information of 3 fields of the third film frame are recorded along the second track turn; and so forth. In this manner, the information of a plurality of film frames are recorded over and with respect to all track turns.\nFor this reason, in the case of obtaining a still picture by reproducing only the same track turn a plurality of times, the information of two different film frames are repeatedly reproduced. Consequently, in the case where the movements of the image of the information contents are rapid, the resulting still picture flickers and becomes indistinct."}
{"text": "In most geographic areas around the world, vehicles are used by persons as a mode of transportation from one geographic location to another. Currently, there are many types of vehicles used for transportation. One type of vehicle is an automobile that includes a chassis supported by four or more wheels (e.g., a car or a truck). Another type of vehicle is a motorized two-wheeled vehicle, which may be in the form of a motorcycle or a scooter. All of these vehicles allow a person to travel a greater distance in a shorter period of time than the person could travel by foot.\nIn many cases, these vehicles are purchased from a third-party provider, sometimes referred as a “dealer”, who may operate as an agent for a manufacturer. For instance, an automobile manufacturer may have hundreds or thousands of dealers at different locations around the world. Depending on the locations of these dealers, the dealer may have to compete against dealers selling other types of automobiles as well as dealers selling the same automobiles in a neighboring geographic region (e.g., neighboring city, neighboring county, etc.).\nWith the increased popularity of the Internet over the last decade, a greater amount of purchasing decisions are being made after reviewing materials on a dealer's website. For instance, a potential customer may access a website of a dealer to review the dealer's vehicle inventory. More specifically, potential customers may review the images for certain vehicles of interest (e.g., exterior and interior images of the vehicle) to determine if the dealer has a particular vehicle of interest and possibly which dealer to visit if she or he wants to inspect the vehicle before purchasing.\nIn some cases, these vehicle images are “lot” images, namely digital photographs of the actual vehicle that are taken by the dealer or by a professional photographer hired by the dealer. For these images, the exterior and/or interior of the vehicle are captured along with the actual area surrounding the vehicle during the photoshoot (e.g., parking lot, parking garage, etc.). The area surrounding the vehicle offers little assistance in differentiating a dealer from its competitors and the use of lot images greatly increases the dealer's overall marketing costs. Also, when the lot images are taken by the dealer, in some cases, these images may include artifacts (e.g., window glare, shadows, less than optimal resolution, etc.), which impede conveyance of the actual features of the vehicle.\nIn other cases, the vehicle images may be generic representations (stock image) for the specific vehicle make/model/trim, and thus, the images may not accurately represent the specific vehicle of interest. As an example, the automobile of interest may include a particular factory-ordered option, such as a navigation system with a compact disc (CD) player for example. However, the stock images and lot images may merely illustrate the interior color and a generic interior layout, not each and every factory-ordered option desired included on the specific (actual) vehicle such as the CD player option with the navigation system.\nThe lack of accurate representations of a specific vehicle of interest along with customized dealer background and/or foreground information may adversely affect a potential customer's decision to purchase the vehicle from that dealer."}
{"text": "Piezoelectric devices, such as piezoelectric quartz filters, piezoelectric quartz resonators and the like; typically include a piece of piezoelectric material mounted to a substrate. In quartz devices, the quartz element of necessity has thin metallic electrodes attached to it through which electrical signals are coupled into and out of the piezoelectric quartz material. Common problems with piezoelectric devices are adequately isolating the piezoelectric device from mechanical shock and dealing with thermal expansion coefficient mismatches between the piezoelectric material and the substrate material.\nQuite often, the piezoelectric devices, such as a piezoelectric quartz material and the substrate have mismatched thermal expansion coefficients. This mismatch can cause mechanical stresses to be induced in the quartz as time goes by during the life of such devices, as the quartz and substrate expand and contract over temperature variations. Further, mechanical shock transferred to the quartz through its mounting structure can increase mechanical stresses that in addition to the thermal stress, adversely affect the output frequency and accuracy of these devices.\nVarious attempts over the years have been developed in an attempt to compliantly mount piezoelectric quartz devices to a substrate. For purposes of this application, a compliant mount for a piezoelectric device, is a mounting device, apparatus or other mounting means that attempts to reduce or minimize mechanical stresses on the piezoelectric quartz element. Some prior art compliant mounting devices have used thin foil tabs that act as spring-type mounting structures that attempt to isolate the quartz element from its substrate. Other types of compliant mounting structures have attempted to use substrate materials having thermal expansion coefficients which more closely match the thermal expansion coefficient of the quartz material itself.\nMost, if not all of the prior art compliant mounting schemes are difficult to use because of the small physical dimensions that modern piezoelectric quartz elements have. Using bent foil tabs alone, for example, to compliantly mount a small silver of quartz onto a substrate is not a structure that lends itself to economic mass production of quartz crystal devices.\nA prior art overtone crystal mount structure is shown in FIG. 1. The crystal mount structure 10 includes a substrate 12, conductive pads 14 attached to an external portion of the substrate for electrical connection to the structure 10, a pair of C-shaped mounts 16, a piezoelectric element 18 including a left wraparound electrode 20 and a right wraparound electrode 22 coupled to the C-shaped mounts by an upper adhesive 24. A lower adhesive 26 securely couples the C-shaped mounts 16 to the conductive pads 14. The adhesives 24 and 26 are epoxy.\nWhen an AC voltage is applied across the pads 14, the piezoelectric element 18 vibrates acoustically at a certain displacement. The resonant frequency of the piezoelectric element 18 drifts about a nominal frequency, in part due to changes of temperature, stresses, exposure to mechanical shock, microcracks in the quartz and the like.\nThere is a need for an improved method of mounting a piezoelectric device with a substrate, to: (i) minimize the mechanical stresses induced due to the thermal expansion mismatches between the two; (ii) provide a mechanically sufficient coupling such that the device can withstand mechanical shocks; and (iii) provide a method of crystal attachment which is adaptable to mass production.\nAccordingly, a low cost, readily-manufacturable, compliant mount for a piezoelectric device would be an improvement over the prior art. A method by which quartz devices can be easily and readily attached to a substrate and which isolates the quartz element from mechanical stresses would be an improvement over the art."}
{"text": "1. Field of the Invention\nThis invention relates to the design and preparation of semiconductor materials having skutterudite type crystal lattice structures with voids or cavities that may be selectively filled to enhance various thermoelectric properties of the semiconductor materials.\n2. Description of the Related Art\nThe basic theory and operation of thermoelectric devices has been developed for many years. Presently available thermoelectric devices used for cooling typically include an array of thermocouples which operate in accordance with the Peltier effect. Thermoelectric devices may also be used for heating, power generation and temperature sensing.\nThermoelectric devices, for cooling applications, may be described as essentially small heat pumps which follow the laws of thermodynamics in the same manner as mechanical heat pumps, refrigerators, or any other apparatus used to transfer heat energy. A principal difference is that thermoelectric devices function with solid state electrical components (thermoelectric elements or thermocouples) as compared to more traditional mechanical/fluid heating and cooling components. The efficiency of a thermoelectric device is generally limited to its associated Carnot cycle efficiency reduced by a factor which is dependent upon the thermoelectric figure of merit (ZT) of materials used in fabrication of the associated thermoelectric elements. Materials used to fabricate other components such as electrical connections, hot plates and cold plates may also affect the overall efficiency of the resulting thermoelectric device.\nThe thermoelectric figure of merit (ZT) is a dimensionless measure of the effectiveness of a thermoelectric device and is related to material properties by the following equation:\nZT=S2\"sgr\"T/xcexaxe2x80x83xe2x80x83(1)\nwhere S, \"sgr\", xcexa, and T are the Seebeck coefficient, electrical conductivity, thermal conductivity and absolute temperature, respectively. The Seebeck coefficient (S) is a measure of how readily the respective charge carriers (electrons or holes) can transfer energy as they move through a thermoelectric element which is subjected to a temperature gradient. The type of carrier (electron or hole) is a function of the dopants in the semiconductor materials selected to form each thermoelectric element.\nThe electrical properties (sometimes referred to as power factor, electrical characteristics, electronic properties, or electronic characteristics) associated with materials used to form thermoelectric elements may be represented by S2\"sgr\". Many of the materials which are used to form thermoelectric elements may be generally described as semiconductor compounds or semiconductor materials. Examples of such materials will be discussed later in more detail.\nThe thermoelectric figure of merit (ZT) is also related to the strength of interactions between the carriers and vibrational modes of the crystal lattice structure (phonons) and available carrier energy states. Both the crystal lattice structure and the carrier energy states are a function of the materials selected to form each thermoelectric device. As a result, thermal conductivity (xcexa) is a function of both an electronic component (xcexae) primarily associated with the respective carriers and a lattice component (xcexag) primarily associated with the respective crystal lattice structure and propagation of phonons through the respective crystal lattice structure. In the most general sense, thermal conductivity may be stated by the equation:\nxcexa=xcexae+xcexagxe2x80x83xe2x80x83(2)\nThe thermoelectric figure of merit (ZT) may also be stated by the equation:                     ZT        =                                            S              2                        ⁢            T                                ρ            ⁢                          xe2x80x83                        ⁢            κ                                              (        3        )            \nxcfx81=electrical resistivity\n\"sgr\"=electrical conductivity       electrical    ⁢          xe2x80x83        ⁢    conductivity    =                    1                  electrical          ⁢                      xe2x80x83                    ⁢          resistivity                    ⁢              xe2x80x83            ⁢      or      ⁢              xe2x80x83            ⁢      σ        =          1      ρ      \nThermoelectric materials for cooling and power generation applications such as alloys of Bi2Te3, PbTe and BiSb were developed thirty to forty years ago. More recently, semiconductor alloys such as SiGe have been used in the fabrication of thermoelectric power generation devices. Commercially available thermoelectric materials are generally limited to use in a temperature range between 200xc2x0 K and 1300xc2x0 K with a maximum ZT value of approximately one. The efficiency of such thermoelectric devices remains relatively low at approximately five to eight percent (5-8%) energy conversion efficiency. For the temperature range of xe2x88x92100xc2x0 C. to 1000xc2x0 C., maximum ZT for many state of the art thermoelectric materials also remains limited to values of approximately 1, except for Texe2x80x94Agxe2x80x94Gexe2x80x94Sb alloys (TAGS) which may achieve a ZT of 1.2 to 1.4 in a relatively narrow temperature range. Recently developed materials such as Si80Ge20 alloys used in thermoelectric generators to power spacecrafts for deep space missions have an average thermoelectric figure of merit (ZT) of approximately equal to 0.5 from 300xc2x0 C. to 1,000xc2x0 C.\nMany crystalline materials with low thermal conductivity do not have good electrical conductivity and many crystalline materials with good electrical conductivity often have relatively high values of thermal conductivity. For example, many binary semiconductor compounds which have skutterudite type crystal lattice structures have relatively good electrical properties. However, the value of thermal conductivity associated with the skutterudite type crystal lattice structure of such semiconductor compounds is generally relatively large which often results in a thermoelectric figure of merit (ZT) which is less than desired.\nResearch and development has previously been conducted on fabricating thermoelectric devices with thermoelectric elements formed from materials having skutterudite type crystal lattice structures. Examples of such developments are shown in U.S. Pat. No. 5,610,366 entitled High Performance Thermoelectric Materials and Methods of Preparation issued Mar. 11, 1997. Patent application Ser. No. 08/101,901 filed Aug. 3, 1993, entitled A Semiconductor Apparatus Utilizing Gradient Freeze and Liquid-Solid Techniques, now U.S. Pat. No. 5,769,943, and patent application Ser. No. 08/412,700 filed Mar. 29, 1995, entitled Advanced Thermoelectric Materials with Enhanced Crystal Lattice Structure and Methods of Preparation, now U.S. Pat. No. 5,747,728, show additional examples of skutterudite type crystal lattice structures which may be used to fabricate thermoelectric devices.\nOne reference presents X-ray analysis of synthesized samples with the major phase being partially filled phosphorous based skutterudite type crystal lattice structures. See, xe2x80x9cSynthesis and Crystal Structure of a New Series of Ternary Phosphides in the System TR-Co-P (TR: Rare Earth)xe2x80x9d by S. Zemni, D. Tranqui, P. Chaudouet, R. Madar, and J. P. Senateur, Journal Solid State Chemistry, 65, 1 (1986).\nIn accordance with teachings of the present invention, the design and preparation of semiconductor materials used in fabrication of thermoelectric devices have been substantially improved to provide significantly enhanced operating efficiencies. The present invention provides the ability to obtain increased efficiency from a thermoelectric device having one or more thermoelectric elements fabricated from materials with skutterudite type crystal lattice structures modified in accordance with the teachings of the present invention to optimize selected thermoelectric characteristics of the resulting thermoelectric device. A significant reduction in thermal conductivity (xcexa) may be achieved by filling a portion of the voids associated with skutterudite type crystal lattice structures as compared to materials having skutterudite type crystal lattice structures with either essentially no filling of the associated voids or approximately one hundred percent (100%) filling of the associated voids. By selecting atoms and/or molecules in accordance with teachings of the present invention to fill a portion of such void spaces, phonon propagation through the resulting skutterudite type crystal lattice structure may be significantly affected to reduced thermal conductivity (xcexa) while at the same time minimizing any reduction in electrical properties (S2\"sgr\") of the associated materials. Examples of semiconductor materials and compounds which may be satisfactory for use with the present invention include, but are not limited to, CoP3, CoAs3, CoSb3, RhP3, RhAs3, RhSb3, IrP3, IrAs3 and IrSb3 and alloys of these compounds.\nIn accordance with one aspect of the present invention, a selected portion of the voids or cavities typically associated with a skutterudite type crystal lattice structure may be filled to optimize the reduction in thermal conductivity (xcexa) of the associated semiconductor material while at the same time minimizing any reduction in electrical properties which results in maximizing the thermoelectric figure of merit (ZT) for the resulting thermoelectric device. Semiconductor materials having a skutterudite crystal lattice structure with a portion of the associated voids filled in accordance with teachings of the present invention have demonstrated an order of magnitude decrease in thermal conductivity (xcexa) measured at room temperature and a much larger decrease at lower temperatures. Thermoelectric devices with thermoelectric elements fabricated from materials with skutterudite type crystal lattice structures partially filled in accordance with teachings of the present invention may have substantially enhanced thermoelectric operating characteristics and improved efficiencies as compared to previous thermoelectric devices.\nIn accordance with another aspect of the present invention, both N-type and P-type semiconductor materials with skutterudite type crystal lattice structures may have a portion of their associated void spaces filled. Thermoelectric elements fabricated from such semiconductor materials may be used for manufacturing thermoelectric devices with substantially enhanced operating characteristics and improved efficiencies as compared to previous thermoelectric devices.\nTechnical advantages of the present invention include the use of semiconductor materials having skutterudite type crystal lattice structures in which a portion of the voids associated with such crystal structures have been filled with atoms and/or molecules selected in accordance with teachings of the present invention. Fabricating a thermoelectric device from such semiconductor materials may substantially enhance the associated operating efficiency. Such thermoelectric devices may be used for cooling, heating, electrical power generation and temperature sensing.\nVarious embodiments of the present invention have been reduced to practice and selected tests have been performed to prove the principles and concepts associated with the present invention. The xe2x80x9cdisorderxe2x80x9d resulting from atoms, selected in accordance with teachings of the present invention, to fill voids or cavities in skutterudite type crystal lattice structures associated with various semiconductor compounds have a substantial influence on reducing thermal conductivity (xcexa) of the resulting lattice structure. Filling only a portion of the voids or cavities in accordance with teachings of the present invention may result in less scattering of the carriers when subjected to both a temperature gradient and a difference in DC voltage as compared to skutterudite type crystal lattice structures with approximately one hundred percent (100%) filling of such voids or cavities. Some of these embodiments and selected test results will be discussed later in more detail.\nDepending upon the desired thermoelectric device, various semiconductor materials may be formed in accordance with the teachings of the present invention include, but are not limited to, the following described general formulas. One variation of the present invention may be used to form a semiconductor compound with the general formula\nX41By1A3xX212xe2x88x923x\nwith y selected to be in the range of greater than zero and less than one to fill the selected portions of normally void subcells to optimize the thermoelectric figure of merit. Another variation of the present invention is a semiconductor compound having a partially filled skutterudite type crystal lattice structure with the general formula:\nDyExX14xe2x88x92xX212\nwith D selected from the group consisting of La, Ce, Pr, Sm, Eu, Gd, Tb, Dy, Ho, Er, Tm, Yb, Lu, K, Th, Pa, U, Sn, Pb, Bi, Br and I; E selected from the group consisting of Co, Rh and Ir; X1 selected from the group consisting of Ni, Pd, Ru, Os, Pt and Fe; X2 selected from the group consisting of Sb, P, As, Bi, Sn, In, Ga, Ge, Si or a combination thereof; and, y selected to be in the range of greater than zero and less than one to fill the selected portion of normally void subcells to optimize the thermoelectric figure of merit. Another embodiment of the present invention includes the following variation to form a semiconductor compound having the general formula:\nX14YbyA3xX212xe2x88x923x\nwith a generally stoichiometric ratio of elements X1 and X2 which normally form a binary semiconductor compound having the skutterudite crystal lattice structure; A selected from the group consisting of Bi, Sn, In, Ga, Ge, and Si; and y selected to be in the range of greater than zero and less than one to fill the portion of normally void subcells to optimize the thermoelectric figure of merit. Lastly, the present invention may utilize the following variation to form a semiconductor compound having the general formula:\nYbyExX14xe2x88x92xX212\nwith E selected from the group consisting of Co, Rh and Ir; X1 selected from the group consisting of Sn, Ge, In, Ga and Si; X2 selected from the group consisting of Sb, P, As, Bi, Sn, In, Ga, Ge and Si, or a combination thereof; and y selected to be in the range of greater than zero and less than one to fill the selected portion of normally void subcells to optimize the thermoelectric figure of merit."}
{"text": "A plasma display panel (hereafter briefly referred to as a “panel”) is a typical AC type surface discharge panel in which many discharge cells are formed between a front plate and a rear plate that are disposed facing each other. In the front plate, multiple display electrode pairs, each including a pair of a scan electrode and a sustain electrode, are formed in parallel on a front glass substrate. A dielectric layer and a protective layer are formed covering these display electrode pairs. In the rear plate, multiple data electrodes are formed in parallel on a rear glass substrate, and a dielectric layer covers these data electrodes. Then, multiple barrier ribs are formed in parallel to data electrodes. A phosphor layer is formed on the surface of the dielectric layer and the side face of the barrier ribs. Then, the front plate and the rear plate are disposed facing each other such that the display electrode pairs and the data electrodes are disposed orthogonal to each other, and sealed. Discharge gas, for example, containing 5% xenon in partial pressure ratio, is filled in an internal discharge space. A discharge cell is formed at an area where the display electrode pair faces the data electrode. In a panel as configured above, an ultraviolet ray is generated by gas discharge in each discharge cell. This ultraviolet ray excites each phosphor of red (R), green (G), and blue (B) to emit light for color display.\nIn general, a subfield method is adopted as a panel-driving method. More specifically, one field period is divided into multiple subfields. Grayscale display is achieved by combinations of subfields to emit light.\nEach subfield includes an initializing period, address period, and sustain period. In the initializing period, an initializing discharge occurs so as to form a wall charge needed for a subsequent address operation on each electrode. At the same time, priming particles (a detonator for discharge=Excited particles) are generated so as to reliably generate address discharge. In the address period, an address pulse voltage is selectively applied to discharge cells to be displayed. This generates address discharge and forms the wall charge (this operation is hereafter also referred to as “address”). In the sustain period, a sustain pulse voltage is applied alternately to the display electrode pair including the scan electrode and the sustain electrode, so as to generate a sustain discharge in the discharge cells where the address discharge has occurred. This makes phosphor layers of corresponding discharge cells emit light, and thus an image is displayed.\nIn the above operations, heat is generated in the discharge cells in proportion to the number of discharges. Accordingly, a temperature of the panel itself increases by this heat. In addition, a brighter display image requires more number of discharges. A brighter display image thus results in a higher panel temperature. Furthermore, it is generally known that a discharge characteristic changes depending on the discharge cell temperature in this type of panels. Accordingly, too high panel temperature causes unstable discharge. This risks degradation in the image display quality.\nTherefore, diversifying methods have been proposed to prevent degradation in the image display quality that may be caused depending on the panel temperature.\nFor example, one method disclosed is to attach the panel to a chassis with a heat-conducting sheet made of silicone rubber in between. In this method, the heat generated from the panel is efficiently transferred to the chassis so as to prevent a temperature rise in the panel. (For example, refer to Patent Document 1.)\nAnother method proposed is to detect the panel temperature and apply diversifying corrections depending on detected temperature. For example, one of plasma display devices disclosed is provided with a panel temperature detector for detecting the panel temperature, and an address pulse period is changed depending on temperature information received from the panel temperature detector so as to stabilize the address operation. (For example, refer to Patent Document 2).\nStill another method disclosed is to detect an average picture level (APL) of an input video signal and the panel temperature, and control a picture level of a display image depending on the detected APL and detected temperature. The panel temperature is controlled by controlling power consumption in the panel. (For example, refer to Patent Document 3.)\nHowever, since distribution of the panel temperature is not even over a panel area, the entire display area hardly shows uniform temperature. In addition, the panel temperature greatly changes depending on display images. On the other hand, reduction of the number of components configuring the plasma display device has been strongly demanded. It is thus difficult to increase the number of thermal sensors to be installed.\nAccordingly, accurate detection of the panel temperature is difficult in a plasma display device in which the number of thermal sensors that can be installed is limited. Still more, a recent trend of larger panel and larger plasma display device makes accurate detection of the panel temperature more difficult.\nFurthermore, in recent years, a brighter display is studied in order to improve the image display quality. In a plasma display device that is configured to increase luminance of light emission by increasing the discharge current, the panel temperature increases further easily.\nAccordingly, it is becoming more difficult to maintain an appropriate panel temperature in the recent trend of a larger panel for plasma display devices and higher luminance.\nHowever, the control of panel temperature to an appropriate level is becoming more important in order to display high-grade images. Therefore, a simple structure for keeping an appropriate panel temperature in plasma display devices for larger panel and brighter image is strongly demanded.    Patent Document 1: Japanese Patent Unexamined Publication No. H10-254372    Patent Document 2: Japanese Patent Unexamined Publication No. 2004-61702    Patent Document 3: Japanese Patent Unexamined Publication No. 2000-305514"}
{"text": "The present invention relates to monitoring systems for inventory levels, and, in particular, a communications based system for monitoring and maintaining inventory levels in agricultural feed bins.\nMany industries, including the agricultural, concrete, and other bulk material handling operations, require ongoing observation and determination of the contents to provide for continuity of supply through timely reordering to prevent service interruptions. For example, large scale confinement operations for livestock store feed in numerous large feed bins with varying contents proximate the livestock facility. Moreover, a single operation may have multiple facilities, each with multiple bins. In as much as the needs of each facility differ based on species, age and diet regimen, there is little redundancy among the plurality of bins, requiring for monitoring of each to insure efficient continuity of operations. Furthermore, each feedstock may have a designated supplier having varying minimum shipment quantities and delivery schedules. Accordingly, it is imperative that the facility operators make regular determinations of inventory status for each bin on a frequent basis.\nAt present, even at large facilities such determinations are extremely difficult and generally prone to error. A typical feed bin is a domed cylindrical metal container having a conical top hatch for the delivery of the feed and a conical bottom material discharge chute.\nTypically, the hatch provides the sole visual access for determining quantities therein. Accordingly, the facility personnel much scale each bin, and visually or with graduated measuring sticks, measure or estimate the quantity therein, record the relevant information, and upon return to the facility offices make determinations on which bins require resupply from designated suppliers in a time frame that obviates the possibility of a food shortage. Alternatively, the personnel might tap the bin and try to correlate the remaining quantity based on audible response. Both of these manual approaches are prone to estimation errors and missed readings that can lead to feed being ordered prematurely and belatedly. Such ordering mistakes, particularly for large operations, are very costly.\nVarious container measuring systems have been proposed for agricultural feed containers. None, however, are entirely satisfactory for the varying types of seed material that have to be handled within a facility, and none provide material assistance in maintaining proper levels throughout the facility to avoid shortages, premature purchasing and other costly incidents. For example, U.S. Pat. No. 5,847,567 to Kielb et al. discloses a microwave system for determining the distance between the top level of a liquid or solid in a tank. In addition to being costly, the system relies on a relatively flat surface for accurate measurement, typically not present in bottom discharge feed bins. The system is also highly prone to error at the bottom conical portion, whereat substantial material volume is carried. U.S. Pat. No. 4,065,967 to Beeston discloses a silo level indicating system using temperature sensors to detect the presence of material at varying levels in a bin. The system relies on temperature variations between a reference sensor at the top of the bin and the contained material. At extreme temperatures, hot and cold, the difference may not be sufficient to enable accurate determinations. U.S. Pat. No. 5,002,102 to Hosel discloses a light detector system for determining the fill level of fiber material in a container. The system also relies on a relatively flat upper surface for measuring the fill level, and like the sonic systems is extremely prone to error at the lower levels of the container. U.S. Pat. No. 5,164,555 to Brenton discloses a rotating paddle that is used to sense the level of material in a container to prevent overfill, rather than determine overall volume of the material remaining. U.S. Pat. No. 4,047,434 at Marsh discloses a sensor system mounted on the side walls of a silo that relies on the conductivity of the material for actuating the sensors for determining the fill level. While the silage contains sufficient moisture for establishing the requisite conductivity, the materials typically stored in feed bins are generally poor electrical conductors, thus erroneous readings are encountered.\nThe above patents are also primarily directed to providing on-site information regarding the fill level by varying visual displays. Thus, while overcoming the need for manual determinations of fill level, the information must be retrieved, dated and reorder procedures conducted elsewhere.\nIn view of the foregoing, it would be desirable to have an automated system for monitoring the level in each bin, and based on current information determine the reorder mechanism for ensuring timeliness and economy.\nThe present invention provides an affirmative measurement of feed levels in a bin by a vertically disposed series of pressure sensors carried on a flexible detector strap releasably mounted at the center of the containers. The sensors include radially projecting blades that deflect and actuate in response to the downward flow of material being discharged. The ability of the blades to dynamically rotate in the presence of movement allows for a more uniform detection of feed at the sensed level. The detector strap is easily mounted in conventional feed bins without requiring structural modifications or additions. The signals from the detector strap are routed to a local monitoring site permitting on-site determination of fill level. The signals from plural bins may also be routed to a central monitoring site. The local monitoring site is preferably provided with a programmable logic controller for periodically recording the level of the feed bin. The fill data is remotely retrievable for use in determining resupply needs. The system may also beneficially incorporate a modem communications link with a remote service monitoring site. The service monitoring site periodically retrieves fill data, determines fill level, reorder requirements and timing, and places an order with a designated supplier to ensure resupply in a time frame insuring feed continuity.\nAccordingly, it is an object of the present invention to provide a container level detecting system that is easily incorporated into existing containers without structural modification.\nAnother object is to provide a container level monitoring system that is effective for varying materials contained in agricultural feed bins.\nA further object is to provide a container level monitoring system that periodically records fill levels in container for retrieval by an off site service provider that determines procedures for reordering."}
{"text": "Farmed animal reproduction has grown into a multi-million dollar business. For years, in the horse racing field, racing champions have been selected to stud mares in hopes of providing an offspring that will similarly be a champion racer.\nIn the avian industry the breeding of avian such as turkeys has become its own industry. In particular, as the turkey industry has become more modernized not only have the turkey moved indoor to confinement facilities, a greater emphasis on efficiency within turkey operations has been emphasized. This is true among all facets of turkey production including reproduction.\nIn the avian field and in particular in the turkey industry artificial insemination is becoming more commonplace in the industry. In particular, semen is collected and then inseminated or delivered with a syringe or plastic straw into the oviduct. This process is expensive and has many drawbacks that have held back more widespread use of this technique. First, once semen is collected its spermatozoal viability can start to decrease after only an hour. Multiple solutions to this issue all have drawbacks, in particular, cooling the sperm to between 2-5° C. in a storage facility and then thawing for insemination can lead to increasing the viability for several hours, though this remains a short period of time. Thus often a semen extender must be added to the semen in addition to being cooled. The cooling storage facilities and extender are expensive and can decrease fertilization.\nIn addition, sperm production in avian is known to be effected by photoperiods. In particular, sperm production increases during spring time, or when light increases are sensed by the avian. Meanwhile, sperm production decreases during the fall time when photoperiods are decreasing over time."}
{"text": "Disclosed herein are curable (meth)acrylate compositions and, more specifically ultraviolet (UV) radiation curable (meth)acrylate compositions. The compositions are suitable for optical articles and particularly for light management films.\nIn backlight computer displays or other display systems, optical films are commonly used to direct light. For example, in backlight displays, light management films use prismatic structures (often referred to as microstructure) to direct light along a viewing axis (i.e., an axis substantially normal to the display). Directing the light enhances the brightness of the display viewed by a user and allows the system to consume less power in creating a desired level of on-axis illumination. Films for turning or directing light can also be used in a wide range of other optical designs, such as for projection displays, traffic signals, and illuminated signs.\nCompositions used to form light management films to direct light desirably have the ability to replicate the microstructure needed to provide the light directing capability upon cure. It is furthermore desirable for the glass transition temperature (Tg) of the cured composition to be high enough for shape retention during storage and use. It is also desirable for light management films made from the cured composition to exhibit high brightness. Finally, the composition used to make light management film advantageously provides a cured composition having a high refractive index (RI). While a variety of materials are presently available for use in light management films, there remains a continuing need for still further improvement in the materials used to make them, particularly materials that upon curing possess the combined attributes desired to satisfy the increasingly exacting requirements for light management film applications.\nThe above-described needs are alleviated by a curable composition comprising a multifunctional (meth)acrylate; a substituted or unsubstituted arylether (meth)acrylate monomer; a brominated aromatic (meth)acrylate monomer; and a polymerization initiator.\nOther embodiments, including a method of preparing a curable composition, a cured composition comprising the reaction product of the curable composition, and articles comprising the cured composition, are described below."}
{"text": "1. Field of the Invention\nThis invention relates to a multi-layered-type liquid crystal display device with the use of a twisted nematic display process that can attain an excellent color display.\n2. Description of Prior Art\nLiquid crystal display devices are now being used in clocks and electronic calculating machines, in displays of computers terminal and word processor, in televisions, and in a variety of other uses in many fields. Recently, there has been an extremely large demand for liquid crystal display devices because of the changes to multicolor and full-color displays, which are already being made use of in the fields of graphic display and image display. Color display, that has been widely put into practical use, is attained by a liquid crystal cell with color filter layers. The liquid crystal cell functions as a light-switcher, and produces various colors. The main kind of display mode is a twisted nematic display mode attained by a liquid crystal cell in which the liquid crystal is twisted at 90.degree.. (The \"twisted nematic\" is referred to hereinafter as \"TN\".) However, when the TN display device is driven at a high duty ratio, the contrast of the image decreases. On the other hand, a supertwisted birefringence effect (which is referred to hereinafter as \"SBE\") process has been suggested, which produces TN oriented liquid crystal which is twisted at angles of approximately 180.degree. to 270.degree. larger than 90.degree.. With the SBE process, the curve in the vicinity of the threshold value increases steeply, and even when the duty ratio increases, it is possible to obtain a high contrast ratio. However, because birefringence effects of liquid crystals are used, the dependence of the display characteristics on the wavelength of light is theoretically higher than with the TN display in which the liquid crystal is twisted at 90.degree.. Thus, it is very difficult to adapt it for use in a full-color display."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the hybrid variety designated CH147656, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include greater yield, better stalks, better roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen comes to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop hybrid parent plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants and have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
{"text": "This invention relates to water-based herbicidal preparations containing at least one agrochemical agent from the group consisting of bromoxynil and its derivatives and certain alkoxylated fatty acid esters, to the use of such fatty acid esters for increasing the penetration of bromoxynil and its derivatives into plant leaves and to a process for controlling unwanted plant growth.\nIt is known that weeds can be controlled with a variety of herbicides which often have little or no solubility in water. One important representative is bromoxynil (3,5-dibromo-4-hydroxybenzonitrile) and its derivatives. Where herbicidal agents such as these are applied, penetration into the plant is minimal on account of their poor solubility in water, resulting in an inadequate herbicidal effect. Accordingly, agrochemical agents of the type in question are formulated as aqueous dispersions or emulsions so that they can be applied without difficulty, for example by spraying onto the plants, the emulsifier also acting as a wetting agent to improve uptake of the agent into the plant. The greater the penetration of the herbicide, the more effectively and efficiently it can be used. Accordingly, there is a general need for formulations which provide for substantially complete penetration of the herbicidal agent into the plant. DE-OS 29 24 403, for example, describes water-based preparations containing bromoxynil salts andxe2x80x94as formulation aidsxe2x80x94alkali metal salts of alkyl polyglycol ether phosphate partial esters. EP 0 485 207 describes aqueous emulsions containing bromoxynil octanoate and heptanoate and, as emulsifiers, polyalkylene oxide-modified silanes or ethoxylated tall oil amides.\nIt has surprisingly been found that water-based preparations containing bromoxynil or derivatives thereof and, as emulsifiers, certain alkoxylated alcohols provide for very high penetration of the herbicide into the plant.\nThis invention relates to water-based herbicidal preparations containing at least one agrochemical agent from the group consisting of bromoxynil and its derivatives and certain alkoxylated fatty acid esters, to the use of such fatty acid esters for increasing the penetration of bromoxynil and its derivatives into plant leaves and to a process for controlling unwanted plant growth.\nAccordingly, the present invention relates to water-based preparations containing bromoxynil or derivatives thereof and one or more compounds corresponding to general formula (I):\nROxe2x80x94(C2H4O)n(C3H6O)mxe2x80x94Rxe2x80x2xe2x80x83xe2x80x83(I)\nin which RO is an alcohol selected from the group of branched or linear, saturated or unsaturated monohydric alcohols containing 1 to 6 carbon atoms or polyols containing 2 to 12 carbon atoms and 2 to 6 hydroxyl groups and Rxe2x80x2 is hydrogen and/or a group xe2x80x94COxe2x80x94Rxe2x80x3, where Rxe2x80x3 is a branched or linear, saturated or unsaturated alkyl group containing 5 to 29 carbon atoms, n is a number of 1 to 50 and m is 0 or a number of 1 to 10.\nThe preparations according to the invention contain bromoxynil and/or derivatives thereof, preferably bromoxynil alkyl esters or sodium and/or potassium salts of bromoxynil or sulfates and carbonates thereof. Preparations containing bromoxynil octanoate or heptanoate are particularly preferred.\nThe alkoxylated compounds corresponding to formula (I) are known substances which are described, for example, in U.S. Pat. Nos. 2,678,935, 3,539,518, 4,022,808 or GB 1,050,497 of which the disclosures are also part of the present invention.\nThe compounds corresponding to formula (I) may by prepared by any methods known to the expert, for example by esterification of fatty acids with alkoxylated methanol, as described in U.S. Pat. No. 3,539,518. Unfortunately, this process has certain disadvantages. It is a two-stage process, the esterification reaction is very time-consuming and the products are colored by the high reaction temperatures. In addition, fatty acid methylester ethoxylates produced by this process have relatively high OH values after the esterification step which can be problematical for certain applications. Another possibility is the direct reaction of fatty acid esters with alkylene oxides in the presence of transition metal catalysts, as described in U.S. Pat. No. 4,022,808. However, the fatty acid alkyl ester alkoxylates are preferably produced by the heterogeneously catalyzed direct alkoxylation of fatty acid alkyl esters with ethylene oxide and/or propylene oxide on calcined or hydrophobicized hydrotalcites. These synthesis processes are described in detail in WO 90/13533 and WO 91/15441 of which the disclosure is also part of the present invention. The products formed are distinguished by a low OH value, the reaction is carried out in a single stage and light-colored products are obtained. The fatty acid alkyl esters used as starting materials may be obtained from natural oils and fats or may be synthetically produced.\nThe alkoxylated fatty acid esters contain at least 1 mol ethylene oxide groups per mol ester. Compounds of formula (I) containing between 1 and 30 mol ethylene oxide per mol ester are preferred. In addition to the ethylene oxide units, between 1 and 10 propylene oxide groups are also present in the molecule. Other preferred compounds corresponding to formula (I) are those which contain between 1 and 30 mol ethylene oxide per mol ester and 1 to 10 mol propylene oxide groups. For these mixed ethylene oxide/propylene oxide adducts, it is possible to use both compounds reacted with a mixture of ethylene oxide and propylene oxide and compounds reacted with ethylene oxide and propylene oxide in two separate steps. The alkoxides are statistically distributed between the OH groups present, depending on the production process.\nIf compounds corresponding to formula (I) containing polyols as the alcohol component RO are used, the quantitative data for the ethylene or propylene oxide units (indices n and m) are always based on the molecule as a whole. However, the exact distribution of the ethylene or propylene oxide units between the various hydroxyl groups of the polyols is known to comply with a distribution dependent on the synthesis process.\nThe fatty acid esters xe2x80x94COxe2x80x94Rxe2x80x3 contain alkyl groups Rxe2x80x3 with 5 to 29 carbon atoms. Suitable fatty acid components are natural or synthetic fatty acids, more particularly straight-chain, saturated or unsaturated C6-30 fatty acids, including technical mixtures thereof obtainable by lipolysis from animal and vegetable fats and oils, for example from coconut oil, palm kernel oil, soybean oil, sunflower oil, rapeseed oil, cottonseed oil, fish oil, bovine tallow and lard; special examples are caprylic, capric, lauric, lauroleic, myristic, myristoleic, palmitic, palmitoleic, oleic, elaidic, arachic, gadoleic, behenic and erucic acid.\nSuitable alcohol components RO are linear or branched, saturated or unsaturated monohydric alcohols containing 1 to 6 carbon atoms, for example methanol, ethanol, n- and i-propanol, n- and i-butanol, pentanol, hexanol, 2-ethylhexanol and cyclohexanol. Suitable C2-6 polyols are, for example, ethylene glycol, 1,2-propylene glycol, 1,2-butylene glycol, glycerol or trimethylol propane and pentaerythritol.\nBasically, all the hydroxyl groups of the alcohols are substituted by the alkoxides although not all terminal alkoxide groups are capped by ester groups. Accordingly, if polyols, such as glycerol or ethylene glycol, are used as the alcohol component RO, the preparations may contain compounds corresponding to formula (I) obtained by reaction both of the full esters and of the partial esters with alkoxides. However, preferred compounds of formula (I) are those in which all the hydroxyl groups of the alcohols are alkoxylated and all terminal alkoxide groups are capped by ester groups with the formula xe2x80x94COxe2x80x94Rxe2x80x3. Accordingly, in these preferred compounds, Rxe2x80x3 in formula (I) stands exclusively for a branched or linear, saturated or unsaturated alkyl group containing 5 to 29 carbon atoms.\nIn addition, alkoxylated fatty acid esters corresponding to formula (I) of which the fatty acid component is selected from linear, unbranched C6-18 fatty acids and of which the alcohol component is methanol, the esters (I) preferably containing between 1 and 3 mol propylene oxide and between 1 and 6 mol ethylene oxide per mol ester, are preferably used in the preparations according to the invention. Compounds such as these can be obtained, for example, by the above-described reactions of palmitic, stearic, oleic, linoleic or linolenic acid, lauric acid and myristic acid or esters thereof with alkoxides.\nAlso suitable are alkoxylated compounds where the alcohol component is glycerol and the fatty acid component is selected from saturated or unsaturated, branched or unbranched fatty acids containing 18 to 22 carbon atoms and the esters contain between 1 and 3 mol ethylene oxide per mol ester. Compounds of formula (I) in which n is 5, 10 or 30 and m is 0 are particularly preferred. Compounds such as these can be obtained, for example, by reacting glycerol esters of natural fatty acids, such as, for example, palm oil, rapeseed oil, soybean oil or preferably castor oil, with ethylene oxide.\nThe compounds of formula (I) present in the preparations according to the invention are nonionic compounds which may also be characterized by their HLB value (hydrophilic/lipophilic balance according to Griffin\"\"s definition; see Rxc3x6mpp, Lexikon Chemie, 10th Edition 1997, page 1764). Preferred preparations contain compounds of formula (I) with HLB values of 4 to 10 and preferably 5 to 9.\nThe preparations contain at least bromoxynil or derivatives thereof as herbicides although mixtures with various other herbicides may also be used. The preparations according to the invention may contain the herbicide in enriched form, in which case they are formulated as concentrates containing more than 50% by weight to at most 90% by weight of herbicidal agent. However, they may also be present in dilute form. Preferred preparations contain between 0.01 and 5% by weight of bromoxynil or derivatives thereof, based on the weight of the preparation. If the preparations contain other agrochemical agents, they are present in quantities of 0.01 to 10% by weight. The percentage water content of the preparations according to the invention is preferably between 10 and 99.9% by weight. The quantity ratio between the compounds of formula (I) and the herbicidal agents is preferably between 1:1 and 1:100. Particularly preferred preparations are those in which the ratio by weight between the compounds of formula (I) and the herbicidal agents is in the range from 1:10 to 1:80 and more particularly in the range from 1:2 to 1:5.\nBesides bromoxynil and its derivatives and the compounds of formula (I), the water-based preparations according to the invention may contain other typical ingredients and additives. These include solvents, such as ethylene or propylene glycols and C1-6 alcohols, solid carriers, such as lignin, lignin derivatives or clays and other known emulsifiers or dispersants. However, preparations containing only emulsifiers corresponding to formula (I) and no other emulsifiers or dispersants are particularly preferred. Most particularly preferred preparations are those which are free from colloids, such as titanium dioxide, and/or free from solvents, more particularly mineral-oil-based solvents.\nThe preparations according to the invention are stable in storage, even at temperatures above 30xc2x0 C., and can be produced without intensive shearing, for example by manual stirring. The preparations according to the invention are formed without intensive shearing, for example by simple manual stirring. To this end, the compounds of formula (I) may be initially introduced into a mixing vessel, for example in liquid form. The herbicidal agent is then added and the resulting mixture is dispersed in water. If compounds of formula (I) with melting points above room temperature are used, they may be used in molten form. However, compounds corresponding to formula (I) with a melting point below 25xc2x0 C. are preferably used. Alternatively, a mixture of the agrochemical agent in water may also be prepared in a first step and the mixture thus prepared may be subsequently emulsified or dispersed by addition of compounds corresponding to formula (I).\nThe present invention also relates to a process for controlling unwanted plants in which a water-based bromoxynil-containing preparation as described in the foregoing is applied to the leaves of the plants by any method known to the expert in such quantities that the plants die off.\nThe present invention also relates to the use of compounds corresponding to formula (I) for increasing the penetration of bromoxynil or its derivatives into plant leaves."}
{"text": "1. Field of the Invention\nThe present invention relates to reflectors and reflective liquid crystal displays, and more specifically relates to a reflector having reflection characteristics such that it appears brighter when reflection light is observed at a specific viewing angle than when it is observed at other viewing angles, and to a reflective liquid crystal display using the reflector.\n2. Description of the Related Art\nLiquid crystal displays can be generally classified into two types: transmissive liquid crystal displays and reflective liquid crystal displays. In reflective liquid crystal displays, external light is used for illumination and a front light is used for ensuring visibility. Reflective liquid crystal displays are commonly used as display units for electronic devices such as mobile computers, calculators, digital watches, communication equipment, game machines, measuring devices, electronic display boards, etc.\nAn example of a reflective liquid crystal display is shown in FIG. 8. With reference to FIG. 8, a display-side substrate 20 and a reflector-side substrate 10 oppose each other with a liquid crystal layer 30 therebetween. The display-side substrate 20 is transmissive and the reflector-side substrate 10 is reflective. The external surface of the display-side substrate 20 serves as a display surface, and the reflector-side substrate 10 is provided with a reflective layer 12. In this reflective liquid crystal display, external light incident on the display surface passes through the display-side substrate 20 and the liquid crystal layer 30, is reflected by the reflective layer 12 in the reflector-side substrate 10, passes through the liquid crystal layer 30 again, and is emitted from the display surface, thereby making an image visible.\nIn FIG. 8, the reflector-side substrate 10 is formed by laminating a glass substrate 11, the reflective layer 12, an intervening layer 13, a color-filter layer 14, a planarizing layer 15, a transparent electrode layer 16 formed of an Indium Tin Oxide (ITO) film, a Nesa film, etc., and an alignment layer 17, in that order from the bottom. In addition, the display-side substrate 20, which opposes the reflector-side substrate 10 across the liquid crystal layer 30, is formed by laminating an alignment layer 21, an insulating layer 22, a transparent electrode layer 23 formed of an ITO film, a Nesa film, etc., a glass substrate 24, and a light-modulating layer 25 (a polarizing plate, a retardation plate, etc.) in that order from the liquid crystal layer 30.\nIn the above-described liquid crystal display, the color-filter layer 14 in the reflector-side substrate 10 includes red (R), green (G), and blue (B) color films which are sequentially arranged in parallel to each other in a striped pattern, and the transparent electrode layer 16 includes transparent electrodes disposed in parallel to each other in a striped pattern at positions corresponding to the color films. In addition, the transparent electrode layer 23 in the display-side substrate 20 includes transparent electrodes which are arranged in parallel to each other and perpendicularly to the transparent electrodes of the transparent electrode layer 16. Parts of the liquid crystal layer 30 at intersections of the transparent electrodes of the transparent electrode layer 23, which is disposed at the display-side, and the transparent electrodes of the transparent electrode layer 16, which is disposed at the reflector-side, are formed as pixels, each pixel corresponding to one of the colors.\nIn addition, in the above-described liquid crystal display, a front light (not shown) is disposed outside the display-side substrate 20 as required. In such a case, similarly to external light, light emitted from the front light passes through the display-side substrate 20 and the liquid crystal layer 30, is reflected by the reflective layer 12 in the reflector-side substrate 10, passes through the liquid crystal layer 30 again, and is emitted from the display surface.\nThe reflective layer 12 in the reflector-side substrate 10 can be generally classified into the specular-reflection type and the diffuse-reflection type. FIG. 9A shows a specular-reflection type reflective layer 12′ and the reflective surface of this reflective layer 12′ is made flat so that the absolute value of the incidence angle and the absolute value of the emission angle with respect to the normal of the display surface are the same. Accordingly, when the display surface is observed, there are problems in that the brightness of the display surface varies depending on the positional relationship between the light source and the viewpoint and visibility is degraded due to back reflection, that is, reflection of the light source, the observer's face, etc., in the display surface. In order to solve such problems, in a diffuse-reflection type reflective layer 12″ shown in FIGS. 9B and 10, a plurality of small concavities and convexities (concave portions 31 in FIG. 10) are irregularly formed next to each other on the reflective surface of the reflective layer 12″. Thus, in the diffuse-reflection type reflective layer 12″, external light incident at a certain angle is diffusely reflected by the surface of the reflective layer 12″. Accordingly, it is possible to obtain a reflective liquid crystal display having a wide viewing angle in which brightness does not vary even when the viewpoint is moved and back reflection is reduced.\nWith regard to the material of the diffuse-reflection type reflective layer 12″, the shape and distribution of the concavities and convexities, and the method for forming the concavities and convexities, various suggestions have been made from the viewpoints of reflection characteristics and productivity.\nRegarding the method for forming the concavities and convexities, a method is known in which light is radiated on the surface of a plate-shaped resin substrate formed of a photosensitive resin layer, etc. through a pattern mask, and a plurality of small, spherical concave portions 31 are formed next to each other by a development process. In order to obtain a mirror-finished surface, a layer of aluminum, silver, etc. is formed on the surface on which the concave portions are formed by vapor deposition or plating. In addition, another method is also known in which a plurality of small, spherical concave portions 31 are formed next to each other by pressing a punch (stamping tool) having a hemispherical end portion against the surface of a flat substrate such as an aluminum plate, a silver plate, etc.\nThe concave portions 31 are generally formed in a spherical shape whose depth varies in the range of 0.1 μm to 3 μm, and distances between the concave portions 31 are set such that the pitch between the concave portions 31 (distance between the central points of the concave portions 31) varies in the range of 5 to 50 μm.\nAn example of a desk calculator is shown in FIG. 11A, and an example of a mobile computer is shown in FIG. 11B. As shown in FIGS. 11A and 11B, when an observer actually views the display surface of a liquid crystal display, he or she often looks up at the display surface from the lower side thereof. More specifically, the viewpoint Ob of the observer is inclined toward the lower side of the display surface by an angle θ relative to the normal X perpendicular to the display surface.\nOn the other hand, in reflective liquid crystal displays, external light is used for illumination; however, the intensity of the external light is greatly reduced as it passes through the light-modulating layer 25 formed of a polarizing plate, etc., the two transparent electrode layers 16 and 23, the liquid crystal layer 30, the color-filter layer 14, etc., and returns. In addition, when the diffuse-reflection type reflective layer 12″ is used, incident light is widely diffused, so that the display surface appears substantially dark when viewed from the viewpoint Ob. Accordingly, when the intensity of external light is small, the visibility is substantially reduced. In the reflective liquid crystal display of the known art, the shape and arrangement of the concave portions are determined such that variations in brightness caused by the difference in viewing angle are made as small as possible. Thus, there is a problem in that sufficient brightness cannot be obtained when the display surface is observed in a specific viewing-angle range, for example, from the lower side relative to the normal of the display surface. In addition, also in the case in which a front light is used, there are problems in that the intensity of the light is reduced and incident light is diffused as in the case of external light. Accordingly, it has been difficult to ensure sufficient brightness in a specific viewing-angle range without increasing the consumption of electrical power for illumination more than necessary.\nAccordingly, reflective liquid crystal displays in which the display surface appears especially bright when viewed in a specific viewing-angle range and back reflection is suppressed over a wide viewing-angle range, have been required."}
{"text": "The present invention is directed, in general, to integrated circuits and, more specifically, to a circuit and method for providing interconnections among individual integrated circuit (xe2x80x9cICxe2x80x9d) chips in a multi-chip module (xe2x80x9cMCMxe2x80x9d).\nIntegrated circuits (xe2x80x9cICxe2x80x9d) are the tiny xe2x80x9cchips,xe2x80x9d usually less than 0.5xe2x80x3 on a side, of silicon (or other similar material) on which is patterned the transistors and interconnections that make modern electronic systems do what they do; compute, amplify, etc. Many advances in IC fabrication processes have brought about dramatic increases in the number of transistors that can be fabricated on each chip, thereby increasing the capabilities, as well as decreasing the size of ICs and, thus, the size of devices using ICs. Whereas the actual size of a typical IC is much less than the size of a conventional IC package, novel IC packaging designs that can significantly reduce the size of electronic devices are being explored. Moreover, as ICs become faster and more powerful, device packaging becomes a major limitation on system speed.\nConventional IC packages consist of the same basic elements: the IC, a lead frame, wire bonds, and an encapsulant. The lead frame is connected to the IC using a very thin wire bonded to both the chip and the lead frame. The encapsulant or molding, usually made from plastic, forms a package that encloses the IC, wire bonds, and part of the lead frame, thereby protecting the IC from the ambient environment. An electronic system is typically constructed from multiple packaged IC devices, which are electrically-and physically-coupled to a printing circuit board (xe2x80x9cPCBxe2x80x9d) by leads that form a portion of the lead frame and which extend from the IC package; the PCB including metallic traces for interconnecting the multiple ICs. A relatively new approach to packaging is to place more than one IC in the same package; the multiple IC package is called a Multi-Chip Module (xe2x80x9cMCMxe2x80x9d), or a xe2x80x9chybridxe2x80x9d package.\nMCM packages are similar to conventional single-chip package designs. MCM packages, however, house more than one IC by mounting conventional chips on a common substrate, which has metallic paths formed thereon that interconnect the individual chips. A conventional lead frame is connected to terminals on the substrate using a very thin wire, and the substrate and lead frame are then enclosed by an encapsulant to form a protective package.\nThe current objective in electronic systems is toward smaller, lighter, faster, portable systems; e.g., cellular telephones, pagers, notebook computers. The development of MCMs may play an important role in furthering that objective by eliminating a level of packaging for many components, facilitating the integration of multiple analog and digital technologies in a single module, reducing electromagnetic interference (xe2x80x9cEMIxe2x80x9d) problems, and increasing the input/output (xe2x80x9cI/Oxe2x80x9d) capabilities per chip. Furthermore, chip-to-chip wiring within a MCM is cheaper and faster than PCB wiring and reduces the board area needed for a device.\nIn some cases, the capabilities of ICs designed for conventional single chip packages is reduced due to a desire to reduce the package size, which necessarily requires reducing the number of package leads. For example, although an IC may process data internally using a plural-conductor (i.e., xe2x80x9cparallelxe2x80x9d) bus, the data may be serialized such that it can be communicated to another IC through only one lead. In an MCM, however, there are no leads associated with individual ICs; i.e., the ICs are coupled internally via very small metallic paths formed on a common substrate that is coupled to a lead frame. Thus, although a principle advantage of MCMs is the capability to integrate many ICs in one package, that advantage is partially diminished if it is necessary to design ICs specifically for use in MCMs, rather than using ICs designed for conventional single-chip packages.\nTherefore, what is needed in the art are circuits and methods for employing conventional ICs in MCMs. There is a further need in the art for techniques of integrating conventional ICs in MCMs such that the performance of the conventional ICs is enhanced when employed in an MCM.\nTo address the above-discussed deficiencies of the prior art, the present invention provides an MCM and methods of operation and manufacture thereof. The MCM includes: (1) a substrate for supporting a plurality of separate integrated circuit (IC) chips thereon, (2) first and second separate IC chips mounted on the substrate, the first separate IC chip including first and second circuit portions coupled together by at least one signal conductor, and (3) interconnecting means that directly couples at least one signal conductor of the first separate IC chip to the second separate IC chip, the interconnecting means bypassing the second circuit portion of the first separate IC chip.\nThe present invention therefore introduces the broad concept of bypassing circuit portions associated with existing IC chips, when used in an MCM, by providing an interconnecting means that directly couples signal conductors within one IC to another IC within the MCM, rather than using the conventional bonding pads used when an IC is separately-packaged. By directly coupling to signal conductors within an IC, circuit portions of the IC can be selectively-bypassed, which may advantageously increase the overall signal processing speed and/or efficiency of the MCM.\nIn one embodiment of the present invention, the interconnecting means directly couples at least one signal conductor of one IC to a conventional bonding pad of a second IC. In an alternate embodiment, the second IC chip also includes first and second circuit portions coupled together by at least one signal conductor, and the interconnecting means directly couples at least one signal conductor of the first IC chip to at least one signal conductor of the second IC chip, thereby completely bypassing the conventional bonding pads of both the first and second IC chips. In effect, the internal signal conductors of the ICs are directly coupled together to form a trans-IC bus that spans the MCM. The interconnecting means, therefore, may selectively-bypass conventional bonding pads and/or circuit portions associated with conventional and existing IC chips. In this manner, conventional and existing IC chips may be suitably employed in conventional single-IC packages, as well as MCMs, without modification thereto.\nIn one embodiment of the present invention, the circuit portions of one or both ICs which are bypassed may be decoupled from a source of electrical power. While not necessary to the broad scope of the present invention, disabling the power to one or more of the circuit portions decreases overall power consumption by the MCM.\nIn one embodiment of the present invention, a first circuit portion of each of the IC chips is a clock driver circuit for providing a clock signal to a second circuit portion of each of the IC chips, the interconnecting means bypassing the clock driver of one of the IC chips whereby the second circuit portion of that IC chip receives the clock signal from the clock driver of the other IC chip. The operation of multiple IC chips within an MCM from a single clock source helps to ensure the accurate transmission of signals therebetween.\nIn one embodiment, described in detail hereinafter, a first IC chip includes a plural-conductor bus and a circuit portion that includes a multiplexing circuit and an output buffer, the interconnecting means directly coupling to the plural-conductor bus to thereby bypass the multiplexing circuit and output buffer of the first IC chip. In a related embodiment, a second IC chip includes a plural-conductor bus and a circuit portion that includes an input buffer and a demultiplexing circuit, the interconnecting means directly coupling to the plural-conductor bus to thereby bypass the input buffer and the demultiplexing circuit of the second IC chip. In combination, the related embodiments provide a means of bypassing the circuits necessary to serialize data communications between the ICs when the ICs are separately-packaged; the interconnecting means providing a trans-IC plural-conductor bus that spans the MCM, whereby the transmission efficiency of data between the separate ICs may be improved when employed in an MCM.\nThe foregoing has outlined rather broadly the features and technical advantages of the present invention so that those skilled in the art may better understand the detailed description of the invention that follows. Additional features and advantages of the invention will be described hereinafter that form the subject of the claims of the invention. Those skilled in the art should appreciate that they may readily use the conception and the specific embodiment disclosed as a basis for modifying or designing other structures for carrying out the same purposes of the present invention. Those skilled in the art should also realize that such equivalent constructions do not depart from the spirit and scope of the invention in its broadest form."}
{"text": "Current sources which are based on a switched-mode converter topology are commonly known. Current sources of this type include at least one inductive storage element, such as a choke, coupled to a load, and a pulse-width modulator. The pulse-width modulator from a constant supply voltage generates a pulse-width modulated (PWM) voltage which is applied to the inductive storage element. In this type of controller, the mean value of a current through the inductive storage element and, therefore, through the load can be adjusted by varying a duty cycle of the PWM voltage.\nProducing the PWM voltage applied to the inductive storage element may involve switching processes which cyclically connect a voltage source providing the constant supply voltage to the inductive storage element. These switching processes may cause electromagnetic interferences (EMI). This is, in particular, an issue when there is a long supply line between the voltage source and a controller including the inductive storage element, like in automotive applications in which the voltage source is, for example, a battery and in which controllers can be distributed throughout the car with supply lines of several meters between the controllers and the battery.\nThere is, therefore, a need to provide a controlled current source in which EMI problems can be prevented or at least reduced."}
{"text": "1. Field of the Invention\nThis invention relates generally to oscillators and, more specifically, to a fast start-up piezoelectric resonator based oscillator which uses a circuit which will prestress the piezoelectric resonator and which will bias the oscillator into a linear operating region after the piezoelectric resonator has been prestressed.\n2. Description of the Prior Art\nThe start-up time of an oscillator is defined as the time required for the oscillator to reach a steady state. Presently, for most oscillators, the start-up time can be several seconds depending on the crystal frequency and amplifier design of the oscillator. The start-up time may be even longer when the temperature of the device using the oscillator increases. This creates the problem since the start-up times may be fairly long and unpredictable.\nTherefore, a need existed to provide an improved oscillator. The improved oscillator must have faster start-up times than present oscillators. The improved oscillator must further have faster and more predictable start-up times even though operating conditions may fluctuate and change."}
{"text": "1. Field of the Invention\nThe present invention relates generally to surgical suturing apparatus and suturing procedures and, more particularly, to suture tie devices preferably and advantageously made of bioabsorbable materials that are particularly useful in endoscopic surgery, applicators for such suture tie devices, and methods of suturing using such suture tie devices.\n2. Discussion of the Related Art\nSuturing of bodily tissue is a time consuming part of most surgical procedures, including both open surgery and endoscopic or closed surgery. As used throughout this disclosure, the term \"open\" surgery relates to surgery wherein the surgeon gains access to the surgical site by way of a relatively large incision formed in the outer surface of the patient's body, whereas the terms \"endoscopic\" and \"closed\" surgery relate to surgery wherein the surgeon gains access to the surgical site by way of one or more relatively small portals formed in the outer surface of the patient's body through which one or more surgical instruments can be introduced to the surgical site. A variety of instruments, such as endoscopes, forceps, cutters, and applicators and the like, can be introduced through the portals to the surgical site. Commonly performed endoscopic surgical procedures include arthroscopy, laparoscopy (pelviscopy), gastroentroscopy, and laryngobronchoscopy.\nPrior to the development of the subject suture device, suturing had been accomplished through the use of a sharp, curved metal suture needle attached to the end of a length of thread-like suture material. The surgeon or surgical attendant would extend the suture needle and trailing suture material through the tissue to be joined by the suture, after which the suture material would be tied into a knot and manipulated such that the knot could be advanced to the tissue site and adjusted for tension in order to accommodate the particular type of tissue being sutured and to control and account for approximation, occlusion, attachment, and other conditions of the tissue. The ability to control tension is extremely important to the surgeon regardless of the type of surgical procedure being performed; however, knotting of the suture material is time consuming and tedious work, particularly in microsurgery and endoscopic surgery. For example, suturing during microsurgery procedures is necessarily tedious and time consuming due to the reduced-size of the suture needle and suture material, and the concomitant difficulty in manipulating the reduced-size suture needle through the tissue and the tying of a knot in the suture material, all of which is performed by the surgeon while viewing the surgical site through various image magnifying devices. With respect to endoscopic surgery, suturing and tying knots represent an even more time consuming procedure due to the narrow confines of the endoscopic instruments. Accordingly, while endoscopic surgery would be preferred for many surgical procedures, the advantages are often outweighed by the disadvantages caused by the length of time required to complete the endoscopic surgical procedure, which time is greatly extended due to the difficulty and amount of time required for suturing.\nThere have been many attempts to provide devices to take the place of conventional suturing with a suture needle and a length of suture material. However, such prior art devices have essentially consisted of staples, clips, clamps, or other fasteners that do not provide the adjustability in tension that can be obtained by the surgeon incident knotting and advancing a length of suture material to the tissue to be sutured. U.S. Pat. Nos. 3,827,277 to Weston; 4,060,089 to Noiles; 4,490,326 to Beroff et al.; 4,513,746 to Aranyi et al.; 4,532,926 to O'Holla; 4,548,202 to Duncan; 4,573,469 to Golden; 4,590,937 to Deniega; 4,595,007 to Meride; 4,602,634 to Barkley; 4,646,741 to Smith; 4,671,280 to Dorband et al.; 4,719,917 to Barrows et al; and 4,741,337 to Smith et al. are representative of such prior art devices for use in place of conventional suturing. Many of these prior art devices are made of bioabsorbable materials such that the devices are absorbed over time into the bodily tissue and do not have to be removed after the tissue has healed.\nThere exist many compositions useful as bioabsorbable materials as represented by the above patents and by U.S. Pat. Nos. 3,739,773 to Schmitt et al.; 3,797,499 to Schneider; 4,141,087 to Shalaby et al.; 4,300,565 and 4,523,591 to Kaplan et al.; and 4,649,921 to Koelmel et al., which discuss characteristics of various bioabsorbable materials and medical devices desirably manufactured of such materials, such medical devices being of a type designed to be engaged in, embedded in, or otherwise attached to various types of bodily tissue, such as bone, muscle, internal organs, skin and other soft tissue, to remain in place in the tissue until the device is absorbed into the body.\nU.S. Pat. No. 3,570,497 to Lemole discloses a suture device formed of a needle with a piercing point extending from a latch cord carrying notches designed to pass through a latch collar. The latch cord is resilient so as to be curved upon itself to form a suture stitch without requiring tying of knot. However, the latching function does not provide the same degree of tension control as can be obtained from knotting a length of suture material. U.S. Pat. No. 4,548,202 to Duncan discloses use of a similar structure in a tissue fastener device, in that serrations or angled barbs are provided along spaced legs passing through tissue to be engaged by an apertured receiver or a flexible filament mesh. U.S. Pat. No. 3,123,077 to Alcamo discloses a surgical suture device having raised projections, depressions or teeth such as barbs or spicules to snag or penetrate tissue so as to hold a sewn incision or wound.\nEndoscopic surgery is generally preferred over open surgery due to the greatly reduced trauma and wound healing time for the patient, and due to the concomitant cost savings associated with shorter hospital stays and the ability to perform some forms of surgery without general anesthesia and in non-hospital or outpatient surgery sites. Accordingly, there has been much effort spent to develop techniques to facilitate suturing so as to further reduce the total duration of surgical procedures. Alternative suturing techniques that have been proposed have included electric coagulation, mechanical devices such as clips, clamps and staples, and lasers. However, no well-accepted alternative has yet been found, since suturing and tying are essential and vital parts of most surgical procedures due to the advantages discussed above with respect to the known alternatives. That is, to date, the proposed alternatives have had disadvantages which outweigh any benefits they may have conferred, chief amongst the disadvantages being an increased health and safety risk to the patient. Furthermore, these proposed suturing devices and procedures do not provide the surgeon with the advantages of suturing and tying and, as a whole, are not useful in a wide range of procedures to allow expansion of the areas in which endoscopic surgery can be effectively performed. Thus, there is a great need for new and advanced suture devices, particularly ones that are applicable in endoscopic surgery, that afford surgeons all of the advantages of knot tying and suture tensioning, and which can be applied in a time-efficient and effective manner."}
{"text": "1. Field of the Invention\nThis invention generally relates to phase-locked loop timing and delay circuitry and, more particularly, to a system and method for providing a timing signal that can be selected from any phase of a voltage controlled oscillator (VCO) output signal.\n2. Description of the Related Art\nFIG. 1 is a schematic block diagram of a timing circuit using a VCO with a four-phase output (prior art). Conventionally, the four phase outputs are 90 degrees apart. In one aspect, the VCO supplies two differential phase outputs, 180 degrees. One of the phase outputs of the VCO can be used as a system clock or loop acquisition, for example. The phase outputs can also be used for other system timing events that are related to the system clock or VCO output signal. However, if the timing event does not occur at 0, 90, 180, or 270 degrees, with respect to the VCO signal, delay elements must be added to achieve the desired phase relationship. For example (as shown), the 90 degree phase output of the VCO may be directed though logic circuits, such as an inverting buffer. The propagation delay associated with logic circuits results in an additional phase shift, so that a resultant 100 degree phase output is created.\nHowever, the propagation delay associated with logic circuits typically varies with respect to lot and temperature. Thus, the phase of the resultant output signal changes. Even if time invariant delays can be added to a system to generate phase shifts, such delays create a phase shift that is completely dependent upon the VCO frequency. If the VCO is operated at a different frequency, the phase shift associated with the time invariant delay necessarily changes, and the output signal phase shifts.\nIt would be advantageous if a VCO output signal could be used as a phase-selectable timing signal.\nIt would be advantageous if the phase of the above-mentioned timing signal could be adjusted to account for changes in the VCO frequency.\nThe present invention system and method provide a timing signal that is derived from a precisely selected phase of a VCO signal. A multiplexer initially selects one of the VCO outputs, which are differentiated by 90 degrees. Then, the phase of the selected signal is finely adjusted through the use of parallel differential amplifiers. As a result, a timing signal can be generated that varies from approximately zero degrees, to approximately 360 degrees of the VCO output signal.\nAccordingly, a method is provided for controlling the phase of a voltage controlled oscillator output. The method comprises: accepting a plurality of VCO outputs coarsely differentiated by phase; selecting one of the plurality of VCO outputs; finely modifying the phase of the selected VCO output; and, supplying the phase modified VCO output.\nAs mentioned above, it is typical that each of the plurality of VCO outputs be differentiated by 90 degrees. That is, accepting a plurality of VCO outputs coarsely differentiated by phase includes: accepting a first VCO output (I0); accepting a second VCO output (I1), differentiated approximately 90 degrees from the first VCO output; accepting a third VCO output (Q0), differentiated approximately 90 degrees from the second VCO output; and, accepting a fourth VCO output (Q1), differentiated approximately 90 degrees from the third VCO output.\nFinely modifying the phase of the selected VCO output includes modifying the selected VCO output in the range between +45 and xe2x88x9245 degrees. In one aspect, the output is modified in 16 discrete steps.\nAdditional details of the above-described method, and a system for controlling the phase of a VCO output are provided below."}
{"text": "1. Field of the Invention\nThis invention relates to a method for inspecting leakage of a sealed container, and more particularly to an inspection method effectively used for inspecting a leaking place such as a pin hole, a crack, etc. of a sealed container containing foods, pharmaceuticals, chemical industrial products or the like.\n2. Description of the Prior Art\nIn the food industry, there are generally used containers of a type which are sealed after foods are filled therein and then sterilized by heating, or containers of a type which are sealed after foods are filled therein under a bacteria-free condition. In these containers, if there occurs a poor sealing, a pin hole or the like, air tends to flow into the containers through the pin hole, etc. from outside and a bacterial contamination takes place. As a result, the food contents of the containers are oxidized, changed in quality, or corrupted. In some cases, it turns out to be a significant problem endangering human life.\nTherefore, in the food industry or container manufacturing industry, many attempts have been made in order to develop an improved art for sealing with a hope to provide a perfect sealing of a container, an improved art for inspecting the sealing of a container, and a related instrument, device, apparatus, etc. Representative examples thereof are as follows.\n(1) One of the representative methods for inspecting leakage of a food can using a metallic can is known as a beat inspecting method. According to this method, an objective container is beaten with a metallic bar and the sealing is inspected by checking with the beating sound. However, this beat inspecting method is inapplicable to a compound container of plastic material with aluminum foils laminated thereon.\n(2) As another example, there is a method for inspecting the sealing of a container, in which an objective container is formed with a hole, the peripheral area of the opening portion of the hole is sealed with a rubber material or the like and then, the container is sunk into water, and thereafter, air is sent into the hole in order to raise the internal pressure of the container and finally, generation of an air bubble from the container is checked. This method, however, has the shortcoming in that breakage of the objective container is unavoidable.\n(3) In a retort food, there is employed a method in that a bag-shaped container formed of a laminated film of aluminum foil and plastic, or a plastic film vapor deposited with aluminum is used, and air is removed from the container after the container is filled with foods, or otherwise an inert gas such as a nitrogen gas or a carbon dioxide gas or a mixture of a nitrogen gas and a carbon dioxide gas is charged into the container after air is removed from the container, and then the container is sealed.\nAs a seal inspecting method of a retort food, there is employed a method in that an objective container containing food therein is kept for two to three weeks at a normal temperature of 25.degree. C. to 40.degree. C. and any expansion of the objective container caused by bacterial contamination or any oozing-out of a liquid content filled in the container, which tends to occur when imperfect sealing, pin-hole, etc. are taken place, is inspected by eye. However, any of the above-mentioned methods has the problems in that many days are required for carrying out this inspection and a large stock of the objective containers is required.\n(4) As still another method for inspecting any leakage, there is known a method in that a container having a confirmed perfect sealing and an objective container are accommodated in two vacuum chambers which are perfectly sealed and separated, and thereafter, any change of an internal pressure of the vacuum chambers, which tends to occur when an imperfect sealing of the objective container, a pin-hole, etc. are detected, by using a so-called pressure sensor under a predetermined reduced pressure. However, this method again has the shortcoming in that when moisture is attached to the external part of the container, a correct leakage inspection is difficult to obtain because the internal pressure of the chambers is also readily changed owing to evaporation of the moisture when the pressure is reduced.\n(5) There is also a method for determining any leakage in that a vacuum chamber is filled with water, and after the objective container is dipped in the water, the internal pressure of the chamber is reduced, and finally, any bubble or any oozing-out of the content of the objective container is inspected. However, this method has the problem that an additional procedure for drying the container is required after inspection because the objective container is wet with the water.\nAs described above, the conventional inspecting methods have the various problems in that the objective container is limited, a correct inspection is difficult to obtain, the manner for carrying out the method is complicated and a long period of time is required for carrying out the method, and so on."}
{"text": "The present invention is directed to circuit that provides a bias voltage to a switching power supply. Devices such as computers, lamp ballasts, battery chargers, and displays use a switching power supply. Various types of switching power supplies are available, but each needs a bias voltage to power the control circuit. The bias voltage, typically 12 volts at 10-100 ma, is derived from an input power source, which may be a main at 120 or 220-240 volts, 50-60 hertz.\nThe presently available bias circuits can be improved by reducing their size, cost and power consumption."}
{"text": "1. Field of the Invention\nThe present invention relates to a mobile communication terminal and corresponding method allowing a user to divide a displayed web page into one or more sub-web pages of a size desired by the user of the mobile communication terminal.\n2. Description of the Related Art\nMobile communication devices such as mobile phones or personal digital assistants (PDAs) provide not only basic functions such as voice call and Short Message Service (SMS) functions, but also various additional functions. Further, mobile terminals now include Phase change Random Access Memory (PRAM) chips as core chips and ultra high pixel Complementary Metal-Oxide Semiconductor (CMOS) chips as camera modules.\nIn addition, mobile terminals can also be used to access the Internet. Therefore, a user of the terminal can access and view a plurality of different web pages. However, because a web page to be displayed is generally larger than a screen of a display unit included in the terminal, the user has difficulty in viewing a web page."}
{"text": "1. Technical Field\nThe present invention relates to a spindle motor.\n2. Description of the Related Art\nGenerally, in spindle motors, sleeves which are provided around rotating shafts rotatably support the rotating shafts, thus ensuring a high level of operating characteristics. Because of these characteristics, spindle motors are widely used as drive units of recording media, such as hard disk drives, optical disk drives, etc., which need high speed rotation.\nSuch a spindle motor uses a fluid dynamic bearing which is operated in such a way that a predetermined amount of fluid is injected between a rotating shaft and a sleeve, which coaxially supports the rotating shaft, so as to form dynamic pressure when the rotating shaft rotates, thus making the rotating shaft rotate more smoothly.\nSpecially, after 2000's shafting systems of spindle motors have largely gone away from using ball bearings to the use of dynamic bearings. Compared to existing ball bearings, the dynamic bearings reduce noise, enhance the impact resistance, and have an extended lifetime. For example, in a fluid dynamic shafting system, there is a dynamic pressure gap between a rotating shaft and a sleeve. A dynamic pressure depression is formed in the sleeve or the rotating shaft. When the rotating shaft rotates with respect to the sleeve, fluid is moved by this relative rotation. Pressure generated by the movement of the fluid serves as a bearing. However, there are problems of friction with the fluid and contact between the rotating shaft and the sleeve.\nIn an effort to overcome these problems, a spindle motor in which a spacing depression is formed in a rotating shaft or sleeve to reduce friction attributable to fluid was proposed. However, for example, in the case where the spacing depression is formed in the rotating shaft, edges formed on the rotating shaft by the spacing depression may come into contact with the surface of the sleeve and thus damage the sleeve.\nTherefore, studies into a spindle motor which can be manufactured by inexpensive processes, such as an ECM process, unlike the conventional technique which includes a pressing process using a plastic deformation principle, and which aims to solve the problem of contact between the rotating shaft and the sleeve which is induced by the edges of the spacing depression formed in the rotating shaft have been conducted."}
{"text": "Graphical user interfaces (GUIs) have been employed as computer interfaces for over two decades. GUIs usually include icons that visually provide access to files and features, menus that textually provide access to files and features, and a pointer icon for selecting the individual icons and menu choices. Two decades ago, the total number of features offered by computer software was sufficiently limited such that they could be presented using the available screen space. However, as the numbers and types of features offered by software multiplied and proliferated, the available screen space was gradually consumed. Unfortunately, the available screen space for GUIs eventually became inadequate.\nAccordingly, there is a need for schemes, mechanisms, techniques, etc. that can efficiently and conveniently enable GUIs to enable access to the multiplicity of features offered by today's software."}
{"text": "Summary of the Invention\nThe invention broadly relates to devices for applying disinfectants and similar pharmaceutical preparations to incisions, wounds and the like, and is specifically directed to a rolling swab applicator for the purpose of applying tinctures and ointments, or as a dry swab for use as a pickup, to aid healing and to treat infection.\nSurgical incisions and other tissue wounds must be carefully treated to prevent infection. For example, after surgery surgical incisions must be dressed and redressed periodically to facilitate healing. An important part of wound dressing is applying a disinfectant such as iodine, hydrogen peroxide or sterile normal saline to prevent infection. This is typically accomplished with the use of a swab applicator which may either be dipped in the disinfectant before application, or which may be prepackaged with the solution saturating the swab.\nConventional swab applicators typically include a handle and an applicator portion connected to the handle that is absorbent and capable of holding a quantity of the disinfectant. The disinfectant is applied by direct engagement of the applicator with the incision or wound.\nAs prescribed, the applicator is wiped or dragged directly over the wound to apply the disinfectant. If the applicator head is circular, an alternative approach is to manually roll the applicator head over the wound. This avoids the transfer of the infection from one portion of the wound to the other. However, conventional swab heads are typically small in circumference, non-rotating and fixed to the handle of the device. A person can manually rotate the handle and cause the swab to rotate as well, but the dexterity required to match the surface speed of the head to that of the handle is difficult and cumbersome.\nThe invention is directed to a rolling swab applicator in which the applicator head is capable of being rotated relative to the handle. Three embodiments are disclosed. In each case, the simple controlled pressure of the swab on the wound allows the swab to travel with the length of the wound as the handle is moved laterally. Controlling the pressure on the swab head on the wound determines the amount of solution wetting the wound, from minimal to flooding the wound with solution as the swab rolls.\nIn the first embodiment, a handle includes an applicator support in the form of an axially projecting split hub. The applicator consists of a spool that can be inserted over the hub for rotation thereon. An absorbent material such as cotton or open cell urethane is wound around the outer periphery of the spool to receive the disinfectant.\nThe second embodiment includes an elongated handle that is formed with an elongated central bore. The applicator includes a small elongated rod that is inserted into the bore for rotation relative to the handle. The rod is preferably retained in the bore by one or more detents. One end of the rod projects beyond the handle and carries an applicator head that is circular in cross section and which has absorbent material placed thereon.\nThe third embodiment is similar to the second in that it includes an elongated hollow handle that receives a shaft with a plurality of axially spaced bearings enabling the shaft to rotate relative to the handle. The shaft is held in place within the handle by a plug member that is inserted into one end of the shaft and which bears against the handle end in a retaining fashion. The opposite end of the shaft projects beyond the handle and carries an applicator head. An absorbent material is wound around or otherwise affixed on the applicator head.\nIn all three of the preferred embodiments, the applicator head rotates relative to a stationary handle, enabling a user to grasp the handle and roll the applicator head over the incision or wound. The absorbent material may first be dipped into a disinfecting solution, or the disinfectant may be applied directly to the absorbent material. Alternatively, the absorbent material may be pre-saturated with disinfectant, sterile normal saline, or ointment.\nThe assembled rolling swab applicator may be provided in different sizes and individually packaged, enabling the user to choose the correct size of the applicator based on the size of the incision or wound. In either case, the absorbent material directly engages the wound in a pressing rather than a wiping manner, depositing or irrigating disinfectant or other solutions directly onto the wound.\nUse of the inventive rolling swab applicator is simple, easy and fast, requiring only opening of the package of an applicator of the appropriate size and rolling the disinfectant onto the incision or wound. The applicator is then discarded.\nThe structure and features of the inventive rolling swab applicator will be more fully appreciated from the drawings and specification."}
{"text": "Central processing units (CPUs) can use virtual addresses to access a physical memory, such as a random access memory (RAM), divided into page units (e.g., of several kilobytes each). Virtual addresses, which differ from physical addresses, are addresses of a virtual memory having a continuous virtual space composed of discontinuous physical memory regions. When a CPU uses virtual addresses, therefore, a memory management unit (hereinafter abbreviated to “MMU”) translates virtual addresses to physical addresses to enable the CPU to access physical memory.\nThe MMU fetches (loads) a page table, which is a translation table for translating virtual addresses to physical addresses, from the RAM into a translation look-aside buffer (hereinafter abbreviated to “TLB”) and translates a virtual address for which the CPU has requested access to a physical address using this page table.\nAs shown in the schematic view in FIG. 4, a page table has memory page entries for access to different memory pages of a physical memory 50. Each memory page entry contains the virtual and physical addresses of the corresponding memory page. The necessary memory page entry is fetched from the page table information stored in the physical memory 50 and is stored in a TLB 54 provided in an MMU 52.\nThe memory page entry also contains a memory protection attribute specifying, for example, “read,” “write,” or “execute.” The technique disclosed in PTL 1 is a technique that sets a value in the protection bit of the memory protection attribute of the memory page entry fetched from the page table into the MMU to prevent a runaway task if a global variable used only by one task is accidentally overwritten by another task."}
{"text": "Terminals may be generally classified as mobile/portable terminals or stationary terminals according to their mobility. Mobile terminals may also be classified as handheld terminals or vehicle mounted terminals according to whether or not a user can directly carry the terminal.\nMobile terminals have become increasingly more functional. Examples of such functions include data and voice communications, capturing images and video via a camera, recording audio, playing music files via a speaker system, and displaying images and video on a display. Some mobile terminals include additional functionality which supports game playing, while other terminals are configured as multimedia players. More recently, mobile terminals have been configured to receive broadcast and multicast signals which permit viewing of content such as videos and television programs.\nAs such functions become more diversified, the mobile terminal can support more complicated functions such as capturing images or video, reproducing music or video files, playing games, receiving broadcast signals, and the like. By comprehensively and collectively implementing such functions, the mobile terminal may be embodied in the form of a multimedia player or device.\nAs such functions become more diversified and complicated, people spend more time on mobile terminals and an application processor (AP) provided in such the mobile terminal has to process more information. To enhance the usability and durability of the mobile terminal, there are more increasing needs for a new heat radiation structure configured to radiate the heat generated in the application processor."}
{"text": "Advanced Mezzanine Cards are printed circuit boards (PCBs) that follow a specification of the PCI Industrial Computers Manufacturers Group (PICMG). PICMG AMC connectors may be used in such a specification. There are problems, however, with the card slots of the connectors when pressing the connectors onto the PCB's. For example, the openings of the card slots may deform or dirt and debris may enter the card opening."}
{"text": "This invention relates to a molded carpet assembly suitable for use for example in automobiles, and to the method of producing such a carpet assembly.\nThe carpeting used to cover the floor areas of automobiles is conventionally molded into a nonplanar three dimensional contoured configuration which conforms to the contours of the automobile floor so as to fit properly, for example, over the transmission hump on the floor of the automobile. In order to make the carpeting moldable and shape sustaining, it is conventionally provided with a backing or coating of a thermoplastic polymer composition. The thermoplastic polymer backing may contain substantial amounts of inorganic fillers to improve the sound deadening properties. Additionally, the automotive carpets are normally provided with padding on the rear surface for cushioning, sound deadening, and thermal insulation. The pads are commonly formed from relatively inexpensive fibers such as jute or recycled waste fibers, and they are precut into predetermined desired shapes and bonded to the back-coated surface of the molded carpet. Formerly, the padding was used over the entire rear surface of the carpet, as shown for example in Squier U.S. Pat. No. 3,673,034. More recently, in an effort to reduce weight and cost and to achieve better conformity to the irregular contours of the automobile interiors, the pads have been applied only in selected areas of the carpet where the cushioning, sound deadening and thermal insulation are most needed. Depending upon the particular automobile body style, the molded carpet may require one or several pads, often of varying thickness, at different locations on the back of the molded carpet.\nIn the production of this type of molded carpet assembly, a number of difficulties are encountered. The pads often vary in size and thickness, which results in nonuniformity of the resulting product. This construction requires manual placement and assembly of the pads which, in addition to the labor costs involved, often results in variations in placement of the pads on the carpet backing, which further contributes to a nonuniform product and results in a poor fit when the carpet assembly is later installed in the automobile. Difficulties are also encountered in obtaining good adherence between the fiber waste pads and the carpet backing. As a result, the pads may have a tendency to become detached during subsequent handling or installation.\nIt has also been proposed to employ a foam cushioning layer instead of fiber pads in moldable automobile carpet structures. Moldable composite carpet structures of this type are disclosed, for example, in U.S. Pat. Nos. 4,016,318; 4,078,100; 4,186,230; 4,230,755; and German OS No. 29 24 197. Typically, they contain both a foam layer and a moldable thermoplastic polymer layer on the back of the carpet. The presence of the foam layer makes molding of the carpet difficult, since the thermal insulating property of the foam interferes with the necessary heating and cooling of the moldable thermoplastic layer during the molding operation. Also, excessive heating may damage the foam layer. In these constructions, the foam pad extends over the entire rear surface of the carpet assembly, and no provision is made for providing the pads in selected areas only of the carpet or for variations in pad thickness, as is needed for many automotive body styles as noted earlier. The presence of the foam pad throughout the rear surface also restricts the shapes and contours which can be imparted to the carpet by molding.\nWith the foregoing in mind, it is an object of the present invention to overcome the disadvantages associated with the existing known methods of production of a molded automotive carpet assembly."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an artificial cotyloid socket having a domal outer shell and a domal inner shell subtending between them a cavity and, more particularly, to such an artificial cotyloid socket wherein the inner and outer shells are non-detachably secured to each other at their rims.\n2. Description of Related Art\nSpherical cotyloid sockets having inner and outer shells that are non-detachably secured to each other at their rims are especially well suited for cement-free implantation into the acetabulum.\nThe state of the art is represented by European patent document A 444,382. This publication discloses such a cotyloid socket in the form of a \"monobloc\" structure. Unfortunately, this design has the drawback that its outer shell is of constant wall thickness and is in a sandwich form. These features restrict elasticity. Suspension between the outer and inner shells must be controlled by an additionally installed spring system and, consequently, is not self-supporting."}
{"text": "Heterogeneous immunoassays require a means be provided to separate a labeled binding reagent from an unlabeled binding reagent. Frequently, a surface is provided to which is bound a specific binding ligand or receptor. Various surfaces have been used, such as latex beads, which can be filtered; tubes or wells, usually plastic, which also serve as the container for the assay mixture; magnetic particles which can be separated in a magnetic field gradient; insoluble polymers which are separated by centrifugation or are used as the stationary phase of a chromatograph; bibulous materials such as cellulose or glass paper through which reagents can be filtered or transferred by capillary action.\nU.S. Pat. No. 4,659,678 describes a method for carrying out an immunoassay in which an antibody labeled with a hapten forms a complex with the analyte in a liquid medium, The complex binds to a solid support to which is bound antibodies to the hapten. Detection of the bound sample is measured using a second antibody to the analyte labeled with reagents such as radioactive iodine, a fluorescer or an enzyme.\nU.S. Pat. No. 4,298,685 describes an assay method using antigens covalently linked to an enzyme, a biotinylated antibody and an avidin coated surface. A competitive immunoassay system is described which is suitable for diagnostic assays.\nEuropean Patent Application 0 201 079 describes a sandwich immunoassay having biotin attached to one antibody, a second antibody attached to a detectable label and a binding substance for biotin attached to a solid support.\nU.S. Pat. No. 4,228,237 describes a method for determining the presence of a ligand in a liquid medium which utilizes enzyme labeled avidin and a biotin labeled reagent. The resulting enzyme activity is related to the quantity of ligand present in the sample.\nUnited Kingdom Application 2 084 317 describes a competitive immunoassay for an antigen which utilizes a solid surface coated with antibody to a hapten or avidin, an antigen-hapten or antigen-biotin conjugate, a soluble antibody to the antigen and an enzyme labeled antibody to the antibody.\nU.S. Pat. No. 4,780,423 describes a heterogeneous assay using controlled pore glass particles. The controlled pore glass particles are used in a fluorescent immunoassay as the support for the specific binding partner capable of binding to a ligand. As used in the invention, the glass particles bind a complex of interest, the detection of which is achieved by use of a fluorescent probe. Measurement of fluorescence is carried out in the presence of the glass particles.\nThe use of solid particles, such as magnetic particles or glass beads, to serve as the support for an immunologic assay is known. An example of such assays include the use of magnetic particles as the solid support in a fluorometric immunoassay as described in U.S. Pat. No. 4,777,145. The use of avidin-coated glass beads in immunoaffinity chromatography and a method for preparing such avidin-coated beads is described by Babashak J. V. and T. M. Phillips, J. of Chromatography 444:21 (1988).\nThe present invention provides a means to carry out various heterogeneous immunoassays using the improved method of the invention so as to achieve high capacity, rapid binding and convenient washing of the stationary phase of the heterogeneous immunoassay without centrifugation or conventional filtration."}
{"text": "1. Field of the Invention\nThis invention relates to poly(ethylene glycol)-poly(orthoester), poly(ethylene glycol)-poly(orthoester)-poly(ethylene glycol), and poly(orthoester)-poly(ethylene glycol)-poly(orthoester) block copolymers.\n2. Background to the Invention\n1. Micellar System for Tumor Targeting\nOne of the major problems in treating cancer is the difficulty of achieving a sufficient concentration of an anticancer agent in the tumor. This is due to the toxicity, sometimes extreme, of such agents which severely limits the amounts that can be used. However, a major discovery in cancer chemotherapy has been the so-called EPR (enhanced permeation and retention) effect. The EPR effect is based on the observation that tumor vasculature, being newly formed vasculature, has an incompletely formed epithelium and is much more permeable than established older vasculature which is essentially impermeable to large molecules. Further, lymphatic drainage in tumors is very poor thus facilitating retention of anticancer agents delivered to the tumor.\nThe EPR effect can be used in cancer targeting by using delivery systems containing anticancer drugs that are too large to permeate normal vasculature, but which are small enough to permeate tumor vasculature, and two approaches have been developed. In one approach, a water-soluble polymer is used that contains an anticancer drug chemically bound to the polymer via a hydrolytically labile linkage. Such drug-polymer constructs are injected intravenously and accumulate in the tumors, where they are internalized by the cells via endocytosis and released in the lysosomal compartment of the cell via enzymatic cleavage of the labile bond attaching the drug to the polymer. Two disadvantages of this approach are that, first, nondegradable, water-soluble polymers have been used, and this requires a tedious fractionation of the polymer to assure that the molecular weight of the polymer is below the renal excretion threshold, and, second, the drug must be chemically attached to the polymer, which in effect creates a new drug entity with consequent regulatory hurdles that must be overcome. The use of polymer conjugates in cancer diagnosis and treatment is discussed in R. Duncan et al., \"The role of polymer conjugates in the diagnosis and treatment of cancer\", S.T.P. Pharma Sciences, 6(4), 237-263 (1996), and an example of an alginate -bioactive agent conjugate is given in Al-Shamkhani et al., U.S. Pat. No. 5,622,718.\nAn alternate approach has been described. In this approach, an AB or ABA block copolymer is prepared where the B-block is hydrophobic and the A-block is hydrophilic. When such a material is placed in water, it will self-assemble into micelles with a hydrophobic tore and a hydrophilic shell surrounding the core. Such micelles have a diameter of about 100 nm, which is large enough that when they are injected intravenously, the micelles can not leave the normal vasculature, but they are small enough to leave the vasculature within tumors. Further, a 100 nm diameter is too small to be recognized by the reticuloendothelial system, thus enhancing micelle lifetime within the blood stream. Additionally, when the hydrophilic block is poly(ethylene glycol), further enhancement of circulation time is noted, as has been observed with \"stealth\" liposomes. The use of block copolymer micelles is reviewed in G. S. Kwon et al., \"Block copolymer micelles as long-circulating drug delivery vehicles\", Adv. Drug Delivery Rev., 16, 295-309 (1995).\nSakurai et al., U.S. Pat. Nos. 5,412,072 and 5,693,751, and Yokoyama et al., U.S. Pat. Nos. 5,449,513 and 5,510,103, describe block copolymers useful as micellar delivery systems where the hydrophilic block is poly(ethylene glycol) and the hydrophobic blocks are various derivatives of poly(aspartic acid), poly(glutamic acid) and polylysine. U.S. Pat. Nos. 5,412,072 and 5,693,751 describe an approach where drugs have been chemically attached to the hydrophobic segment; while U.S. Pat. Nos. 5,449,513 and 5,510,103 describe an approach where hydrophobic drugs have been physically entrapped within the hydrophobic portion of the micelle. This latter approach is clearly preferable because no chemical modification of the drug is necessary.\n2. Bioerodible Block Copolymer Matrix for Controlled Drug Delivery\nIn AB, ABA, or BAB block copolymers comprising a hydrophilic A block and a hydrophobic B block, the A and B blocks are incompatible and on a microscopic scale will phase-separate. This phase separation imparts unique and useful thermal properties to the material.\nThere is considerable prior art in the development of block copolymers comprised of poly(ethylene glycol) and bioerodible hydrophobic segments such as poly(L-lactic acid), poly(L-lactic-co-glycolic acid) copolymers and poly(.epsilon.-caprolactone), and discussion of their use as drug delivery agents. For example, see W. N. E. Wolthuis et al., \"Synthesis and characterization of poly(ethylene glycol) poly-L-lactide block copolymers\", Third European Symposium on Controlled Drug Delivery, 271-276 (1994), L. Youxin et al., \"Synthesis and properties of biodegradable ABA triblock copolymers . . . \", J. Controlled Release, 27, 247-257 (1993), and Bezwada et al., U.S. Pat. No. 5,133,739.\nPoly(orthoesters) are known as potential vehicles for sustained release drug delivery. See, for example, J. Heller, \"Poly (Ortho Esters)\", Adv. Polymer Sci., 107, 41-92 (1993), and references cited therein, Heller et al., U.S. Pat. Nos. 4,304,767, 4,946,931, and 4,957,998, and PCT International Publication No. WO97/25366.\nThe disclosures of these and other documents referred to in this application are incorporated herein by reference.\nHowever, no block copolymer systems have been described where the hydrophobic, bioerodible segment is a poly(orthoester)."}
{"text": "The following outlines OpenFlow, as one example of SDNs (Software Designed Networks) that can dynamically set or change a configuration, function, or performance of a network base on a software. OpenFlow switch includes a flow table including a plurality of flow entries. Each flow entry includes a match field to be matched with header field information of a received packet, a counter field including statistics information such as the number of received packets and the number of received bytes, and an action field with zero or more actions that dictate how the switch handles a received packet whose header matches with the match field. Upon reception of a packet, the OpenFlow switch retrieves the flow table thereof using header field information of the packet. In the case of miss-hit (non-match), the OpenFlow switch forwards the packet to OpenFlow controller over a secure channel.\nBased upon information on a source and a destination of the packet specified in the header thereof, the OpenFlow controller computes a path for the packet from network topology information. Based upon the path, the OpenFlow controller generates and sets a flow entry for each of OpenFlow switches on the path. On reception of following packets, each having a header matching with a match field of the flow entry set by the OpenFlow controller, each of the OpenFlow switches on the path forwards the packets to a next node, for example, as prescribed in the action field of the flow entry. Regarding details of OpenFlow, reference may be made to NPL 6.\nThere have been extensive research and development efforts on secure networking, especially in SDN to reduce network vulnerabilities to attacks.\nIn NPL (Non Patent Literature) 1, there is proposed ROSEMARY controller, which implements a network application containment and resilience strategy based around the notion of spawning applications independently within micro-NOS (network operating system).\nIn NPL 2, there is proposed PermOF coping with potential trust issue on OpenFlow applications, in which abuse of trust could lead to various types of attacks impacting an entire network. In PermOF, isolation of control flow and data is established between a controller and applications. OpenFlow applications are isolated from Controller kernel, that is, OpenFlow applications cannot call kernel procedures or directly refer to a kernel memory. There is provided Access control layer between OpenFlow applications and OS (Operating System). The layer is controlled by the Controller kernel, so that undesirable interaction between OpenFlow applications and OS would be cut off.\nIn NPL 3, there is proposed AVANT-GUARD, a new framework to advance security and resilience of OpenFlow networks with greater involvement from the data-plane layer. Connection migration enables data plane to shield control plane from saturation attacks. Actuating triggers automatically insert flow rules when the network is under attack.\nIn NPL 4, there is disclosed a secure hypervisor, called TrustVisor, to provide a safe execution environment for security-sensitive code modules without trusting OS or application that invokes the code module. TrustVisor is so designed as to protect security-sensitive code and data from malware, even on untrusted commodity platforms.    [NPL 1] Seungwon Shin et al., “Rosemary: A Robust, Secure, and High-performance Network Operating System”, CCS'14    [NPL 2] Xitao Wen et al., “Towards a secure controller platform for openflow applications”, HotSDN '13    [NPL 3] Seungwon Shin et al., “AVANT-GUARD: scalable and vigilant switch flow management in software-defined networks”, CCS'13    [NPL 4] Jonathan M. McCune et al., “TrustVisor: Efficient TCB Reduction and Attestation”, Mar. 9, 2009, (revised Mar. 10, 2010) CMU-CyLab-09-003, CyLab Carnegie Mellon University Pittsburgh, Pa. 15213    [NPL 5] Jan Medved, Reinaldo Penno, OpenDaylight Update, November 2013, Internet <URL: http://www.ietf.org/proceedings/88/slides/slides-88-netconf-6.pdf>    [NPL 6] “Openflow Switch Specification” Version 1.0.0. (Wire Protocol 0x01), Internet<URL: http://www.openflowswitch.org/documents/openflow-spec-v1.0.0.pdf"}
{"text": "The last decade has witnessed the widespread adaptation of handheld wireless voice communication devices (e.g. Cellular Telephones). These devices have revolutionized personal communication by allowing telephone access from anywhere within reach of wireless network infrastructure (e.g., cellular networks, communication satellites, or other infrastructure of other wireless networks adapted for voice communications). Wireless voice communication technology has effected society at the root level of interpersonal relations, to with, people now expect to be reachable and to be able reach others from anywhere.\nIn as much as the use handheld wireless voice communication devices is not restricted to homes and offices, such devices will often be used in environments where there is considerable ambient noise. Examples of such environments include busy urban settings, inside moving vehicles, and on factory floors. Ambient noise in an environment can degrade the intelligibility of received voice audio and thereby interfere with users' ability to communicate.\nIt would be desirable to provide an improved method and apparatus to increase the intelligibility of speech emitted into noisy environments.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity."}
{"text": "Conventionally, a shopping cart is designed to be nestable into a preceding, similarly constructed cart and to enable a following, similarly constructed cart to be nested into the shopping cart. Commonly, in retail stores and in parking areas near retail stores, long lines of nested shopping carts may be formed, which must be moved by store personnel.\nTypically, a rear caster comprises a horn, which unlike the horn of a front caster does not swivel and which mounts an axle, and a wheel, which is journalled on the axle, via a bearing, and which has a tread. Moving of a line of nested shopping carts stresses rear casters of the nested carts, tends to cause the horns of said casters to bend, and tends to cause the bearings and treads of said casters to wear excessively, particularly as attempts are made to turn the line of nested shopping carts.\nIn U.S. Pat. No. 6,923,456 B2, the disclosure of which is incorporated by reference herein, a shopping cart being nestable similarly and having a rear caster-lifting capability is disclosed. The caster-lifting capability reduces wear and damage to rear casters of the shopping cart."}
{"text": "This invention generally relates to foliage mowers which have height-adjustable mower decks that are operated by foot pedal action. More particularly, the invention relates to mowers which are of the type having cutter decks which can be selectively raised and lowered between, for example, one or more heights suitable for transport and one or more heights suitable for foliage cutting. The invention facilitates movement among such positions and securement of the cutter deck at the selected position.\nFoliage cutting vehicles, mowers and implements have long been known which have the ability to raise and lower cutter decks in order to accommodate different cutting heights and/or to facilitate transport. Rose U.S. Pat. No. 4,120,136 shows a tractor having a mower deck which can move vertically through the use of a foot pedal which achieves height adjustment by engagement with one of a plurality of notches adjacent the foot pedal. The position of the pedal in the notches adjusts the position of a lift linkage and thus of the cutter deck. Trefz et al U.S. Pat Nos. 5,138,825 and 5,351,467 describe riding lawn mowers having pedal operated height adjustment mechanisms. Notched or threaded mechanisms remote from the foot pedal cooperate in the raising and lowering function. Gordon U.S. Pat. No. 5,946,893 employs a deck lift mechanism wherein a foot pedal arm engages a height pin for giving a desired cutting height, while requiring a linkage which connects a foot pedal arm to a forward pivot rod, as well as independent suspension structures. Busboom et al U.S. Pat. No. 5,816,033 shows a hand-operated control lever and linkage arrangement. An arcuate frame is adjacent the manual control lever, these two components being laterally offset from each other.\nFoot-operated lift mechanisms have different operational advantages from hand-operated lift mechanisms. In many circumstances, a foot-operated mechanism can more readily allow the user to provide the leverage and force needed to achieve vertical movement of a unit having substantial weight, such as a cutter deck and blade assembly of an industrial or professional quality mower. Heretofore, foot-operated lift mechanisms for mower decks, such as those noted above, utilize complicated adjustment mechanisms or linkage mechanisms. These can increase the cost and/or reduce the reliability of the adjustable height cutter decks.\nIn accordance with the present invention, a foliage mower and mower deck lift mechanism therefor are provided. The deck height lifting is practiced by foot pedal depression which directly lifts the cutter deck. Pivot lugs of a height adjustment assembly interconnect by a push-pull rod assembly to insure that the cutter deck is lifted uniformly. Its corners preferably are moved simultaneously, either upwardly or downwardly, to the same extent. This height adjustment assembly further includes a stop mechanism which prevents downward movement of the cutter deck beyond any one of a selectable plurality of deck height levels during the action of cutter deck lowering. The invention is suitable for mid-mounted cutter decks of mowers, self-propelled riding mowers, vehicles or implements. Mid-mount mowers having self-propelled capabilities and very tight turning radius attributes, such as those having a turning radius approximating zero, find special application for this invention.\nIt is accordingly a general object of the present invention to provide an improved foliage mower structure having one or more cutter decks having a height adjustment feature which includes foot pedal operation.\nAnother object of this invention is to provide an improved self-propelled riding foliage mower having a foot-operated deck height adjustment assembly which provides for sure and positive uniform deck movement in the vertical direction.\nAnother object of the present invention is to provide an improved foot-operated cutter deck lift mechanism which affords direct lift control when the foot pedal is depressed.\nAnother object of the present invention is to provide an improved self-propelled foliage mower and deck lift mechanism which incorporate a stop mechanism to maintain an selected desired deck height.\nThese and other objects, features and advantageous of the present invention will be apparent from and clearly understood through a consideration of the following detailed description."}
{"text": "1. Technology Field\nThis disclosure relates to a tray and, more particularly, to a tray for bearing different storage units.\n2. Description of the Related Art\nStorage units such as hard disks, hard disk cards, memory cards, and PCB modules are commonly born by trays. Accordingly, the storage units can be loaded at apparatus or systems via carriers. Generally speaking, different carriers have to be provided according to different types and specifications of the storage units. In other words, one specific carrier commonly applies to only one storage unit of a specific specification or in a specific type."}
{"text": "1. Field of the Invention\nThe present invention relates to the mass production of thin, flexible electrical low-cost circuits. More particularly, the process according to the present invention aims at the super-mass production of flexible Radio Frequency [RF] planar resonant circuits for one way use and including inductive and capacitive components, as are commercially desired as marker means in Electronic Article Surveillance [EAS] systems, in Electronic Personnel Access Control [EPC] systems, in Electronic Anti-Counterfeit [EAC] systems, and in most modern Object Handling Control [OHC] systems, e.g. electronic document safe sysems, clinical supply systems, electronic flight baggage and ticket control systems, and the like.\nThe invention introduces the computerized nonstop super-mass production of such circuits, and breaks down prior art barriers in view of limits for custom designs of shape, face, and resonant properties of such circuits. Custom designs of circuits are easily layouted onscreen of a modern personal computer. The data obtained thereof are downloaded to the process line. Hopping custom circuit designs with ease is unprecedentedly introduced. Forming tools have been abandoned. As an exclusive longlife production tool electromagnetic wave energy is extensively used. The circuits may be washed and sewn and nicely go one way with a pretty face.\n2. Descripton of the Prior Art\nU.S. Pat. No. 3,913,219 (LICHTBLAU) describes a production process for the make of planar electrical circuits comprising a dielectric sheet and having an inductively reacting conductive spiral path made from an appropriate metal foil and bonded to one side of said dielectric sheet. Said spiral path terminates, at each end, via at least one connection means, into first and second conductive areas being in a mutually aligned opposite and spaced by said dielectric sheet, thus forming a discrete capacitor. Said capacitor and said conductive spiral thus cooperate to form an LC resonant circuit. The production method may be described by putting some light on the major problems encountering therewith, as follows. The production starts up with a trilayer compound material resulting from a difficult collamination process of two aluminum layers of widely differing thickness (typical 8 .mu.m and 50 .mu.m) with a polyethylene layer. The latter must exhibit very close tolerances in view of thickness (normally about of 25 .mu.m +/-5 .mu.m) and dielectric permittivity, and hence has to be produced off-line, in a special apparatus, at extra speed, and thus is expensive. First waste results between its shock cooled high density extrusion and its heated collamination between said aluminum layers, at a considerably low speed. Then, applying known roto-printing techniques, etching resist patterns are printed, in a mutual registration, on both faces of the compound and dried, at high speed. Due to the roto-printing technique, the print is of the fixed formatted type, which means, the circuit patterns repeat at a constant pitch or spacing on the circumference of the printing cylinder. Though printing is performed with a high repetitive register accuracy of the resist structures defining the resonator to be constructed, this repetitive registration is maintained only a very short time from the following reason. When etching, the entire processing web starts up shrinking, owing to strains frozen in the polyethylene during recrystallization in the early moment of its collamination with said aluminum layers aided by heat. These strains get active, once wide areas and spirally turn-to-turn spacings have been removed, and thus cancel any previously achieved registration on the previously nonetched web both in its along and across direction. Further, etching taken alone is a slow process step, since modern rapid etching techniques--due to their inherent poor overetch protection--cannot be utilized, from the following reason. Etching must take place on both faces of the compound simultaneously, meaning, that the etching of aluminum layers having a widely different thickness (e.g. 8 .mu.m and 50 .mu.m) should be finished right at the same time. Overetching and dropped-registration-waste is a result. Shrinking encounters the problem, that the endless processing web representing a plurality of resonant structures cannot be combined, as circuit's face, with an endless cover material preprinted with an individual outline per circuit repeating at a constant resonator pitch or space, since due to lost registration no precise alignment of a resonator and a preprinted circuit face to mate and cover each other can be maintained. Further, since the inconsistent multilayer construction of the fully finished circuit from other reasons can be printed merely asynchronously by low quality rubber-flex techniques, this known process neither yields standard circuits printed with a fixed print-to-resonator registration, nor yields circuits according to hopping an changing custom design specifications, e.g. including a super performance color print. Further problems encounter with that method, since immediately upon etching the processing web, due to its filigree shape, is most complicated to handle, to advance and to position, just in a state in which such should be easy and carried out quickly. Again waste is a result. For an improvement, it is not allowed to thicken the polyethylene layer for stabilization, otherwise the RF performance of resonators would dramatically drop whenever such circuit should be made as small as possible which is, however, a common desire. Further problems encounter with the need, that the two metal layers have to be contacted through said polyethylene layer's normal thickness of of 25 .mu.m. Such thru-contacts proved to be unreliable; \"dead\" circuits use already to appear among good ones when shipped from the factory. Chemical agents as being in worldwide use with the recreation of fashions and garments not only destroy these contacts. These agents dissolve and float off the roto-printable etchant resistive pattern, which is left back surplus in the circuit construction once etched. Hence, circuits of this make may even mediate or cause heavy degradations and deteriorations of suchlike treated fashions in a much later state, provided such circuits are made just for one way use for walking away with customer's buy. Moreover, though said polyethylene layer is considerable costly and its properties are essentially needed merely in a considerable small place of the entire resonator construction, this layer is most inefficiently wasted over the entire circuit area (wasted precision is wasted money). As an overall result, these problems and the restrictions for the finished product concluded thereof qualifies this method as non-applicable for making disposable one way RF marker circuits, which also serve for quite other functions, a variety thereof having been disclosed in more recent patent applications of applicant.\nU.S. Pat. Nos. 4,482,874 (RUBERTUS & TAIT) and 4,598,276 (TAIT) describe a method of producing novel high efficiency marker devices the unique design scheme of which having been taught in U.S. Pat. No. 4,694,283 granted to applicant. This method is dedicated to a complex stamping routine which utilizes a most delicate domino stampling tool. It has a fixed design and requires steady performance control, service and exchange on elapse of its short lifetime. This means, that such production cannot be run continuously 24 hours/365 days a year, at minimum dead time. However, a minor advantage is obtained with this method over the one mentioned above, so far as the pitch or spacing between resonators may be influenced at some degree, however, by encountering complex mechanics along the production line, provided such resonators should be covered with a registered paper face on-line. However, owing to a limited dimensional precision of any cover material to join the endlessly produced resonator web, preprinting such cover material may not be applied without earning all the same lost-registration problems with the attempt to have individual resonator structures superpose in a perfect alignment individual face prints.\nU.S. Pat. No. 4,658,264 (BAKER) also repeats the basic construction scheme of the planar stripline resonator as disclosed in applicant's U.S. Pat. No. 4,694,283. However, it is not taught how to assign to such resonator a precisely registered print on its usable outer face, or how to design such resonator for facilitating its fitting between most modern cover materials, for giving it even a tex-like appearance or the capability of being washed and sewn. Oppositely, this Patent centers merely on the task of how to fold over a substantially filigree and inconsistent substrate for a registered alignment of folded over portions with a precision of alignment sufficient enough, as has been reflcted in applicant's prior publications.\nU.S. Pat. No. 4,369,557 (VANDEBULT) repeats major portions of the LICHTBLAU process, however, with an important distinction so far as the trilayer compound as a production carrier is exchanged for a dual layer one. Thus, etching of a constant thickness metal layer is carried out merely on one face of a thin and dimensionally unstable dielectric layer which, however, again and all the same has to serve for production carrier's function.\nConsequently, all other problems and restrictions of the LICHBTLAU process are prefectly imported and maintained. Hence, that Patent does not teach how to produce one way RF marker circuits offering a modern cover face designed according to whatsoever customer's specifications.\nEuropean Patent Application No. 84104174.2 (SALAMONE, HOOVER & VANDEBULT) discloses a method of fabricating respective circuits by electrolytically depositing and growing up resonator structures. The process applies fixed format silk screen printing electrically conductive patterns on a thin dielectric film, as may be of polyester. Then, this film and patterns are subjected to a very slow electrolytic deposition routine. Since this method, too, does not yield a marketable product rather than mere (bare) resonators, it is also conflicted with the described lost-registration tragedy, so that this process also has to be considered as off-limits for the super mass production of RF marker circuits going both one way and customized, reading with a whatsoever pretty commercial face.\nThe present invention not only accomplishes the huge object of overcoming all of the above problems. It defines a new horizon for making mass marketable electrical one way circuits featuring custom designs to fill catalogs. Therefore, the present invention breaks with prior art methods, for defining a new reference standard of the state of the art, which becomes obvious from an impressive listing of objects of the invention as follows:\nand"}
{"text": "Machine learning algorithms attempt to learn to solve classification problems (e.g. recognizing which handwritten digit has been written from an image) or regression problems (e.g. estimating the price of a financial asset based on various financial variables) by training on a large database of examples where the correct answer is given. They are then typically tested on unknown examples where the correct answer is withheld.\nThere are a variety of existing machine learning algorithms. The currently leading classes of algorithms for classification include those based on neural networks (particularly popular nowadays are so-called “deep belief networks” (DBNs) trained using the back-propagation algorithm), support vector machines (SVMs) and K-nearest neighbor (KNN) algorithms. Versions of all of these algorithms for regression problems also exist, although they are sometimes less natural.\nThe performance of a machine learning algorithm can be assessed in several ways. The most popular approach is to assess based on the average error an algorithm makes on a set of test examples, after very many training examples have been presented. Unfortunately, the best machine learning algorithms, as assessed according to this perspective, are extremely slow to train. They typically need very large training corpuses, and the training time can be as long as several CPU-years, which means that clusters of computers and/or GPU acceleration often need to be employed. Existing state-of-the-art machine learning algorithms also “learn” in a way that seems unnatural compared to human learners—for example, if presented with a training example and its correct answer, they may fail to give the correct answer even if the very same example is used as a test query immediately afterward, and only learn the correct answer after thousands of presentations of the example. Also, they tend to perform very poorly if only a small number of training examples are presented, and they often require that training examples be presented in large batches. Finally, some machine learning algorithms become progressively slower to respond to test queries as they learn from larger training sets."}
{"text": "1. Field of the Invention\nThis invention pertains to the field of recreational and amusement games. More particularly, it pertains to a playing cardbased game for two or more persons that simulates, to a high degree of reality, the esoteric aspects of the game of football.\n2. Description Of The Prior Art\nCard games have existed as far back as the ancient Chinese cultures. There is something both mysterious and exciting about holding a multiplicity of pasteboards, hoping to overpower the opponent with this hand or the next, while relying on skill and the probabilities inherent in all games of chance. Even those whose game-playing skills are not fully developed still seek to play.\nSome card games have their own reward such as the poker game pot or the bridge grand slam. Other card games, either by themselves or played along with special boards and dice, etc., engage in mock combat normally held in other arenas such as war, football games, baseball games, real estate transactions and stock transactions. It is to this type of board game, namely the engagement in mock football games, that this invention is directed.\nFootball games have been simulated with cards for many years. More recently, the game has been simulated by the use of board games containing electronically actuated vibrating means to provide movement of football player figures fixed on the board to progress up and down a simulated playing field. More recently, computer components have been assembled to simulate movement of football players over a simulated field. Scoring is based upon a combination of random selection of numbers or directions coupled with the skill of the player in pressing buttons or actuating switches with sufficient dexterity to overcome a general resistance built into the electronic circuitry. In all of these attempts, however, none has fully simulated control of individual teams using the plays typical to the football game such as running, passing, and kicking, and achieving any realistic degree of simulation with actual conditions encountered in football games. As a result, these simulated games lose their appeal after a period of time, either by virtue of the players adapting to the skill requirements of the switches or other mechanical features needed to be manipulated during the game, or becoming disinterested because of the loss of control of the outcome of the game due to the total random selection of the game elements.\nThose who have once played football still seek the thrill of formulating offensive and defensive football plays in an environment closely simulating the actual football game itself and utilizing their mental skills to overpower the opponent now that their physical skills are no longer sufficient to allow them to personally engage in an actual football game. Accordingly, there is a significant need for a playing card-based simulated football game which allows the players to utilize their mental skills in planning and executing the various offensive and defensive plays in an environment that closely simulates the actual game itself."}
{"text": "The present invention is a bat particularly adapted for use as a practice bat. The bat of the invention is intended for use in practice to develop batting strength, bat speed and hand-eye coordination. In order to accomplish these objectives, as will be described, one of the features of the invention is that the practice bat is somewhat heavier than a conventional baseball bat of equal length and it includes a shaft portion which is substantially thinner than that of a conventional bat.\nFor reference to prior art bats other than coventional baseball bats attention is directed to the Niederberger U.S. Pat. No. 3,179,412, issued Apr. 20, 1965 and the Anderson U.S. Pat. No. 2,964,321, issued Dec. 13, 1960."}
{"text": "1. Field of the Invention\nThe present invention relates to a rotation angle-adjustable rotating device, more particularly to a rotation angle-adjustable rotating device for use in an electrical appliance.\n2. Description of the Related Art\nA conventional rotating device is shown to be employed in an electrically operated fan, and includes coaxially disposed first and second gears 12, 13, a third gear 14 formed with a coupler shaft 15, a connecting plate 16, and a swing arm 17. The first gear 12 is driven by a threaded portion 111 of a drive shaft 11 of the fan, while the second gear 13 is meshed with the third gear 14. One end of the connecting plate 16 is securely mounted on the coupler shaft 15 for co-rotation therewith, and the other end thereof is pivoted to a first end of the swing arm 17. A pushing element 171 is fixed on a second end of the swing arm 17, and extends into an elongated slot 181 formed in a fan housing 18 of the fan in such a manner that rotation of the drive shaft 11 results in synchronous rotation of the coupler shaft 15, which, in turn, results in a back-and-forth swing action of the swing arm 17, thereby driving the fan housing 18 to swing through the pushing element 171.\nOne disadvantage of the conventional rotating device resides in that the rotation angle of the fan housing 18 is unadjustable, thereby inconveniencing the user when he or she desires a larger swinging angle of the fan."}
{"text": "As suggested in U.S. Pat. No. 3,979,273 to Panzera et al, it is well known to sputter deposit both sides of a substrate by rotating the substrate between targets in the vacuum chamber. The substrate is in the shape of a cylinder which therefore requires that the substrate be rotated during the deposition process. There is no disclosure in this patent of a means for depositing on both sides of a substrate at one time without rotation by supporting a substrate through the center of the cathode targets.\nU.S. Pat. No. 3,741,996 to Urbanek et al suggests a cathode/target that is larger than the substrate to be coated. This patent to Urbanek et al is typical of the prior art where the substrate is supported by a platform table which supplies the cooling and power for operation of the sputter deposition. The substrate is not rotated during the deposition process but only one side of the substrate is coated at one time because of the supporting table covering the second side of the substrate. Further, the substrate is not supported through the center of the cathode/target nor is it apparent from this patent how the support could be accomplished.\nU.S. Pat. No. 4,073,262 to Scheffel et al suggests apparatus for the manufacture of magnetic coating of magnetic disks consisting of a means for coating both sides of a disk by spraying the coating material onto both sides of the substrate at one time. The substrate is supported by a ram extending from the center of a cover disk which prevents the sprayed material from escaping into the atmosphere.\nThe cover plates also provide an aerodynamic effect on the spreading of the dispersion over the disks during a spinning operation. The patent does not discuss the apparatus that would be involved in the sputter deposition of a similar substrate. There is no teaching in the patent how a substrate could be supported from the center of a cathode/target to accomplish a sputter deposition of both sides of the substrate without revolving the substrate.\nU.S. Pat. No. 3,189,535 to Webb suggests a dual cathode target system for coating a spherical substrate. The sphere is supported on one end and rotated within for complete coverage. There is no disclosure of a means for supporting the substrate by way of the support from the center of the cathode target for even coating.\nAn IBM Technical Disclosure Bulletin entitled \"Coating Magnetic Storage Disks,\" April 1969 at page 1611 discloses apparatus consisting of two plain parallel cathodes between which a magnetic disk substrate is located affixed to a rotating anode turntable. There is no showing of the apparatus by which the dual coating is attained and it does not appear that the substrate is supported parallel to the two parallel cathodes.\nU.S. Pat. No. 3,878,085 to Corbani; Chapin, The Planar Magnetron, Vacuum Technology Research/Development, January 1974, pages 37-40; and Cormia et al, Magnetically Enhanced High-Rate Sputtering, Proceed. ElectroChemical Soc., September 1974 are of interest with respect to magnetically-enhanced sputtering."}
{"text": "1. Field of the Invention\nThe present invention relates to multilayer electronic components and methods for producing multilayer electronic components, and particularly, to a multilayer electronic component in which external terminal electrodes that are connected to internal conductors are directly formed on surfaces of a ceramic body by plating and a method for producing the multilayer electronic component.\n2. Description of the Related Art\nRecently, the market for compact portable electronic devices has been expanding, including cellular phones, notebook PCs, digital cameras, and digital audio equipment. Such portable electronic devices are becoming smaller while having higher capabilities. Accordingly, smaller sizes and higher capabilities have also been demanded for numerous multilayer ceramic electronic components mounted on the portable electronic devices. For example, there has been a demand for monolithic ceramic capacitors having smaller sizes and higher capacitances.\nOne effective approach to provide monolithic ceramic capacitors having smaller sizes and higher capacitances is to reduce the thickness of ceramic layers. Recently, capacitors including ceramic layers having a thickness of about 3 μm or less have been produced. Although a further reduction in the thickness of the ceramic layers has been attempted, it has been difficult to ensure quality because a reduced layer thickness increases the risk of a short circuit between internal electrodes.\nAnother approach is to increase the effective area of the internal electrodes. However, mass production of monolithic ceramic capacitors requires side margins between the internal electrodes and side surfaces of a ceramic body and end margins between the internal electrodes and end surfaces of the ceramic body to allow for misalignment during lamination and cutting of ceramic green sheets. Therefore, to increase the effective area of the internal electrodes, the area of the ceramic layers must be increased to accommodate predetermined margins. However, there is a limitation on increasing the area of the ceramic layers within predetermined dimensional specifications of the products, and the thickness of external terminal electrodes is a hindrance to increasing the area of the ceramic layers.\nTypically, external terminal electrodes of monolithic ceramic capacitors are formed by applying a conductive paste to the ends of a ceramic body and firing the paste. A common method for applying the conductive paste is to dip the ends of the ceramic body into a paste bath and to remove the ends therefrom. However, in this method, the conductive paste tends to be thickly applied at the central portions of the end surfaces of the ceramic body due to the viscosity of the conductive paste. Thus, the external terminal electrodes are partially thick (specifically, more than about 30 μm thick), and the area of the ceramic layers must be decreased accordingly.\nIn response to this problem, a method in which external terminal electrodes are directly formed by plating has been proposed.\nIn this method, plating films precipitate on exposed portions of internal electrodes at end surfaces of a ceramic body, and as the plating films grow, the exposed portions of the adjacent internal electrodes are connected together. This method enables the formation of thinner and flatter electrode films than the method using a conductive paste (see International Publication No. WO 2007/049456).\nHowever, if the external terminal electrodes are directly formed on the surfaces of the ceramic body by plating, the following problem occurs.\nIf the external terminal electrodes are directly formed by plating, plating growth proceeds in two stages: a primary growth stage in which plating films grow from the exposed portions of the individual internal electrodes primarily in a spreading direction, and a secondary growth stage in which the plating films grown at the primary growth stage combine together in the spreading direction and a plating film grows from the continuous plating film primarily in a thickness direction. However, a problem arises in that it takes a significant amount of time for the separate plating films to grow and combine together in the spreading direction at the primary growth stage because the exposed portions of the internal electrodes are independent and isolated from each other. This increases the time required to form a plating film with a predetermined thickness, thus causing low productivity."}
{"text": "The present invention relates to a pad that provides hypo/hyperthermia properties to a person using the pad.\nIn U.S. Pat. No. 5,336,708, Chen discloses a gelatinous elastomer composite article. These articles, as disclosed by Chen, xe2x80x9cinclude: GMG, MGM, MG1G2M, M1M2G1G2, M2M1G1G2, G1MG1G2, MG1G2M, G1G2M, GM1M2G, G1M1G2M2M1, M1GM2GM3GM4, [sic] ect, where G = gel and M = material The subscript 1, 2, 3, and 4 are different and are represented by n which is a positive number. The material (M) suitable for forming composite articles with the gelatinous elastomer compositions can include foam, plastic fabric, metal, concrete, wood, wire screen, refractory material, glass, synthetic resin, synthetic fibers, and the like. Sandwiches of gel/material . . . are ideal for use as shock absorbers, acoustical isolators, vibration dampers, vibration isolators and wrappers. For example the vibration isolators can be [sic] use under research microscopes, office equipment, tables, and the like to remove background vibrations.xe2x80x9d U.S. Pat. No. 5,336,708, col. 3, lines 35-51. Chen further discloses, xe2x80x9cgenerally the molten gelatinous elastomer composition will adhere sufficiently to certain plastics (e.g., acrylic, ethylene copolymers, nylon, polybutylene, polycarbonate, polystyrene, polyester, polyethylene, polypropylene, styrene copolymers, and the like) provided the temperature of the molten gelatinous elastomer composition is [sic] sufficient high to fuse or nearly fuse with the plastic. In order to obtain sufficient adhesion to glass, ceramics, or certain metals, sufficient temperature is also required (e.g., above 250xc2x0 F. [121xc2x0 C. ]xe2x80x9d U.S. Pat. No. 5,336,708, col. 9, lines 8-18 (brackets added for consistency of temperature comparison).\nElkins in U.S. Pat. No. 4,884,304 describes a bedding system with selective heating and cooling of a person. That system has, from top to bottom, in order: a top mattress cover, a gas envelope and a multiplicity of liquid flow channels. The multiplicity of liquid flow channels is accomplished by a conventional hypo/hyperthermia blanket. The details of this conventional blanket are set forth in this patent. A problem with this system occurs when a person is on the mattress cover. When the person is on that mattress cover, the person has two sides: (1) a xe2x80x9ccontacting sidexe2x80x9d that touches the mattress cover and (2) the xe2x80x9cexposed sidexe2x80x9d that does not touch the mattress cover. The person disperses the gas envelope and only certain portions of the contacting side contact the flow channels. As shown in FIG. 5 of that patent, the shoulders and other peripheral points of the contacting side of the person, such as arms, do not contact the flow channels. Thereby, that bedding system fails to transfer the desired temperature of the flow channels uniformly to all,sections of the contacting side of the person.\nM. Figman in U.S. Pat. No. 3,266,064, and von der Heyde in U.S. Pat. No. 5,887,304 illustrate conventional convective medium mattress system which essentially has a lower xe2x80x9cbox springxe2x80x9d and a mattress made of rubber, foam, or conventional mattress materials that an individual or object lies thereon. In each embodiment, the lower box spring has a cavity that the medium enters and distributes throughout. The medium then escapes from the cavity through apertures of the mattress.\nA problem with these apertures 89 is that they kink 90 when an adult lies 22 thereon, as shown in FIG. 8. Please note that von der Heyde\"\"s system is designed for an infant, not an adult. And an infant is of such low weight that kinking is essentially nonexistent.\nWhen kinking occurs, the medium is prevented from contacting the body. And when the medium does not contact the body, the medium is unable to treat the hypothermia or hyperthermia portions of the patient that contact the mattress, or even cool or heat the portions of the patient that contact the mattress.\nThe present invention solves this problem.\nThe present invention relates to a first conformable material having a three-dimensional shape and a first hypothermia and/or hyperthermia device."}
{"text": "A remote vehicle control system may be used to permit a user control of one or more vehicle devices from a distance. A typical remote vehicle control system, for example, in a more modern vehicle may include a controller connected to a plurality of vehicle devices via a data communications bus, and a remote transmitter providing the controller with remote control signals. The vehicle devices, for example, may operate the engine starter, the door locks, the power windows, or the alarm system.\nAn example of such a system is U.S. Pat. No. 5,719,551 to Flick, which discloses a remote transmitter that can remotely control a number of vehicle devices. The controller in the Flick '551 patent is responsive to the remote transmitter and sends command codes over a data communications bus to the vehicle devices. Similarly, U.S. Pat. Nos. 6,275,147, 6,756,885, 6,756,886, and 6,812,829, to Flick also disclose a controller/transmitter used to remotely control a number of vehicle devices via command codes sent over a communications data bus.\nU.S. Pat. No. 5,583,479 to Hettich et al. also discloses a vehicle controller connected to a number of vehicle devices via a data communications bus and a remote transmitter in communication with the controller. The Hettich et al. patent further discloses that if the alarm system is not deactivated correctly, then the vehicle devices will be impaired or prevented from working properly.\nU.S. Pat. No. 6,232,873 to Dilz et al. discloses a vehicle security system that detects if an original control unit is no longer active. If the security system determines that the original control unit is no longer active, an alternate circuit that remains active in the vehicle activates an electronic immobilization system.\nAlthough conventional remote vehicle control systems operating via the data communications bus have made significant advances in convenience for the user, the overall security may still be an issue. For example, a would-be thief gaining access to the data bus, such as from under the vehicle, may generate rogue commands on the data bus compromising vehicle security. The would-be thief could temporarily connect a rogue controller to the data bus and cause the windows to roll down or the doors to unlock. Once inside the vehicle, the would-be thief could again connect to the data communications bus and start the engine. Of course, if the vehicle had a vehicle security system, the would-be thief could disarm the vehicle security system via the data communications bus."}
{"text": "This invention relates to an apparatus for body surface mapping, for use in diagnosing different conditions of the human myocardium.\nThere are various methods which present the information contained in an ECG to the cardiologist, the most successful of which has been the standard twelve lead ECG. Unfortunately the standard ECG, in many instances, fails to provide an unequivocal diagnosis.\nEven before the conventional twelve lead ECG became a medical standard, electrocardiographic body surface mapping (BSM) was being investigated as a method which increased spatial resolution and thereby increased diagnostic capability. These xe2x80x9cmapsxe2x80x9d were viewed as pictures presenting lines joining areas of the same electrical potential (isopotentials) at a specified instant of time. This is shown in FIG. 1, which is an isopotential map of a normal healthy subject mid-way through the heart\"\"s QRS (depolarization) period. In FIG. 1 the map has been superimposed onto the outline of a human torso and is obtained from multiple ECG electrodes located substantially fully around the torso from the anterior (left hand side of FIG. 1) to the posterior (right hand side of FIG. 1).\nSince each map shows the isopotentials at only a single point in time, to see the whole electrical xe2x80x9cpicturexe2x80x9d requires the viewing of successive maps at successive time instants throughout at least part of the cardiac cycle. Furthermore viewing these maps their pattern and morphology to be considered from one instantaneous map to the next. Thus, for example, FIG. 2 illustrates six body surface isopotential maps showing how map morphology changes with time during the QRS period, the ECGs on the left of FIG. 2 showing the successive points in time during the QRS period that the maps represent. Since the body is cylindrical in shape the maps actually represent the body xe2x80x9cunwrappedxe2x80x9d onto a two dimensional plane (see FIG. 1). It is conceivable, however, that such maps could be produced which do not xe2x80x9cunwrapxe2x80x9d the torso.\nHowever, BSM is difficult to use and as a result much research has been carried out which attempts to bring the benefits of BSM in a fast easy to use form. Methods were introduced of portraying BSM information as mathematical integrations of the individual ECG waveforms. Isointegral maps present the body surface as lines which join areas possessing the same integral values of pre-defined sections of the ECG. Such isointegral maps have shown that there is more information outside the spatial scope of the standard twelve lead which could be used by a clinician to improve patient management, and have proved their ability to provide an accurate diagnosis in instances where the standard twelve lead was equivocal.\nThe use of isointegral maps and discriminant function analysis rose from a need to speed the processing of body surface map information, reduce the data storage requirements of recordings and provide ultimately some form of automated diagnostic system. The mathematical nature of these isointegrals provides a platform which easily lends itself to the use of computers thereby allowing a diagnostic evaluation in minimal time. This diagnosis takes the form of a statistical discriminant function. Statisticians using discriminant function analysis to analyze isointegral maps are presented with the problem of firstly establishing just how many dimensions are involved with solving the problem. This is necessary in order to choose an appropriate discriminant algorithm. However, due both to the integration with respect to time of the ECG signals and the use of a statistical analysis which ignores the nature of the heart\"\"s electrical field, the isointegral maps still fail to provide a diagnosis in many instances.\nIt is an object of the present invention to provide an apparatus for body surface mapping which can provide an improved diagnostic capability and is relatively easy to use.\nAccordingly, the present invention provides an apparatus for body surface mapping, comprising a plurality of electrodes for attachment to spatially separate locations on a human torso, each electrode being capable of detecting the electrical activity associated with a heartbeat and producing a corresponding voltage, means for sampling the voltages on the electrodes a plurality of times during at least part of a cardiac cycle to produce a like plurality of sets of sampled values, and means for calculating and displaying in graphical form the variation with respect to time of at least one characteristic of the sampled values.\nIn a preferred embodiment the characteristic, for example the instantaneous overall maximum or minimum of the sampled values, varies in position with respect to time and the calculating and displaying means calculates and displays the said variation in position. Preferably the characteristic is displayed as projections of the trajectory onto two planes perpendicular to each other and to the plane containing the electrodes.\nThe invention also provides a method of body surface mapping, comprising locating a plurality of electrodes at spatially separate locations on a human torso, each electrode being capable of detecting the electrical activity associated with a heartbeat and producing a corresponding voltage, sampling the voltages on the electrodes a plurality of times during at least part of a cardiac cycle to produce a like plurality of sets of sampled values, and calculating and displaying in graphical form the variation with respect to time of at least one characteristic of the sampled values.\nIn the context of this specification the xe2x80x9ctrajectoryxe2x80x9d of a characteristic means the notional path traced out by that characteristic in a three dimensional space of which two dimensions correspond to mutually perpendicular directions in the plane containing the electrodes and the third dimension is time represented as a third spatial direction perpendicular to the plane containing the electrodes."}
{"text": "The present invention broadly relates to compositions useful in separating metal ions or other target material including, but not limited to, polypeptides, nucleic acids, and endotoxins, from non-target material."}
{"text": "For background reference is made to U.S. Pat. Nos. 5,016,230 and 6,084,504 incorporated herein by reference."}
{"text": "Example embodiments of the present inventive concept relate generally to integrated circuit devices and, more particularly, to embedded refresh controllers and memory devices including the embedded refresh controllers.\nSemiconductor memory devices for storing data may be classified into volatile memory devices and non-volatile memory devices. Volatile memory devices, such as dynamic random access memory (DRAM) devices, may be configured to store data by charging or discharging capacitors in memory cells, and lose the stored data when power is off. Non-volatile memory devices, such as flash memory devices, may maintain stored data even when power is off. Volatile memory devices are widely used as main memories of various apparatuses, and non-volatile memory devices are widely used for storing program code and/or data in various electronic devices, such as computers, mobile devices, etc.\nIn volatile memory devices, cell charges stored in a memory cell may be lost by a leakage current. In addition, when a wordline is transitioned frequently between an active state and a precharged state (i.e., when the wordline has been accessed intensively or frequently), an affected memory cell connected to a wordline that is adjacent to the wordline frequently accessed may lose stored charges. Charges stored in a memory cell may be recharged before data is lost by leakage of cell charges. Such recharge of cell charges is referred to as a refresh operation, and a refresh operation may be performed repeatedly before cell charges are significantly lost."}
{"text": "Among the most effective ways to assign functions to genes is by having genetic variants of the gene in question. Such genetic variants can be of natural origin or generated deliberately by mutagenesis. The effectiveness and use of induced mutagenesis depends a great deal on the mutation efficiency of the reagent and/or the procedure. Various forms or reagents to intentionally induce mutations in plants exist, but chemical mutagenesis by ethyl methanesulfonate (EMS), an inducer of point mutations in DNA, has gained interest. EMS is an alkylating agent that reacts mainly with the guanine (G) base in the DNA to change it to O-6-ethylguanine, which prefers to base pair with adenine (A) during DNA replication instead of its natural partner, the cytosine (C). As a result, an A is introduced at the location in the DNA where there was a C before. Because the new A will naturally hydrogen bond with thymine (T) during further rounds of DNA replication, the original GC base pair is converted into an AT base pair transition when mutagenesis is done with EMS.\nSeed mutagenesis is the method of choice in almost all plant species except maize, where EMS mutagenesis is done via pollen. The main reason for this is that maize is a monoecious plant with separate male and female inflorescences at two separate locations on the plant. The male flowers are present only in the tassel, which forms at the top of the main axis of the plant, and the female flowers are present in the ear, which is a modified lateral shoot on the side of the main axis approximately 5-6 leaves below the tassel node. The maize embryonic cells that eventually morph into the tassel or the ear are already split into separate lineages by the time the maize kernels (seed) mature. When the maize seed is treated with EMS, the mutations that are generated in the ear primordial cells are completely independent from the mutations that are induced in the tassel primordial cells. To get these mutations in homozygous condition, which is absolutely needed to reveal recessive mutations, the M3 generation in maize is of interest because even when M1 plants (derived from EMS treated seed) are self-fertilized to produce the M2 progeny, the ear and tassel mutations are consolidated together in a heterozygous condition. This reproductive peculiarity of maize also causes the original germ line mutations to be amplified enormously, especially via the pollen, imposing serious limitations on the recovery and genetic interpretation of most mutations.\nTo eliminate these problems with seed mutagenesis, as well as to avoid the tremendous amount of work and expense that goes with having to raise an extra generation to recover mutant phenotypes (propagating thousands of plants, especially by self-pollination, is not trivial), a pollen mutagenesis protocol involving EMS was developed. In addition to allowing the revelation of dominant mutants in the M1 generation and recessive mutations in the M2 generation, this pollen mutagenesis protocol generates mutations that are independent of each other. An additional advantage of mutagenesis via pollen is that it allows one to conduct targeted mutant screens (also known as non-complementation screens) to specifically generate a series of mutations in a gene of interest, for example screening for non-complementation of a recessive mutation.\nOne disadvantage of pollen mutagenesis, which has been optimized to work only in maize thus far, is its mutation efficiency. It is about one out of 1,000 gametes treated with EMS, when gauged by functional genetic tests. When tested molecularly, the density of mutations caused by the pollen EMS protocol in maize was found to be about 2 mutations per Mb of DNA. Whatever the reasons, the relatively low mutation frequency of mutagenesis in this protocol acts as a deterrent for many working either in the public or commercial sector. This impediment can be eliminated if the mutation frequency is enhanced.\nThere is therefore an unmet need for a cost-effective, high accuracy method for generating independent mutant alleles."}
{"text": "A method of determining crystal orientation in a wafer is described in the publication \"A latching accelerometer fabricated by the anisotropic etching of (110)-oriented silicon wafers\" by Dino R. Ciarlo in the journal \"Micromech. Microeng.,\" 2, 1992, pages 10-13. According to the method described therein, rectangular mask openings such as circular scale marks are arranged one next to the other. The angle between two adjacent rectangular mask openings is 0.1.degree.. Each rectangular mask opening is 8 .mu.m wide and 3 mm long. During anisotropic etching of the wafer surface to which the etching mask has been applied, undercut zones, which are bounded by (111) planes, develop at the sides of the mask openings. The dimensions of the undercut zone will vary depending on the orientation of the mask opening relative to the crystal orientation which is to be found. The dimensions of a particular undercut zone are determined optically, and the location of the mask opening at the side of which the undercut zone is smallest is utilized for determining the crystal orientation. The dimension of the smallest undercut zone must be larger than the wavelength of the light used for measuring in order for it to be perceived; thus it must be larger than around 0.5 .mu.m. This is ensured according to this method, among other things, by the dimensions selected for the rectangular mask openings. The achievable angle precision in determining crystal orientation using this method depends on the angle between two adjacent mask openings. For reasons related to the manufacture of masks, this angle cannot be made smaller than 0.1.degree.. The best precision achievable in determining crystal orientation is thus 0.05.degree.. The required etching time in the known method is around 30 hours, depending on the length of the rectangular mask openings. While a reduction of the length of the rectangular mask openings would result in a reduction in etching time, it would also affect the dimensions of the undercut zones causing them no longer to be measurable."}
{"text": "In general, amines are organic compounds comprising a functional group that contains a basic nitrogen atom with a lone pair, which can bind a proton to form an ammonium ion. Amines and, particularly chiral amines, are ubiquitous in biology and industry. Amines, particularly chiral amines, are key building blocks in many pharmaceutical, agrochemical and chemical applications. Chiral amines are particularly important in the production of physiologically active compounds. Moreover, in many applications of chiral amines, only one particular optically active form, either the (R) or the (S) enantiomer has the desired physiological activity.\n(Chiral) amines can be produced both by chemical and biocatalytic synthesis routes. Chemical synthesis of optically active chiral amines via a one-step procedure requires high chemo-, regio-, diastereo-, and enantiocontrol. One of the biocatalytic synthesis routes to obtain chiral amines uses the enzyme transaminase (EC 2.6.1.X; also known as aminotransferases). Transaminases are pyridoxal phosphate dependent enzymes and catalyze the transfer of an amine (—NH2) group from an amine donor, for instance an amino acid or a simple amine such as 2-propylamine, to a pro-chiral acceptor ketone, yielding a chiral amine as well as a co-product ketone or alpha-keto acid, in the presence of the cofactor pyridoxal phosphate (PLP) which is continuously regenerated during the reaction (FIG. 1). Transaminases, particularly (R)- and (S)-selective transaminases, have received much attention as suitable catalysts for producing chiral amines, as they allow the direct asymmetric synthesis of (optically active) chiral amines from pro-chiral ketones.\nAlthough the transaminase catalysed synthesis of chiral amines presents a high enantio- and regioselectivity, the transamination reaction is a reversible reaction, often with an unfavourable thermodynamic equilibrium which limits obtaining high chiral amine yields. Accordingly, an amine mixture is obtained comprising the chiral amine and the amine donor, requiring further purification of the chiral amine. In addition, the transaminase catalysed chiral amine synthesis is prone to substrate and product inhibition. In addition, substrate solubility issues may hinder the reaction as well.\nBörner et al. (Organic Process Research & Development, 2015, 19(7), 793-799) describes an enzymatic chiral amine synthesis process setup combined with a selective solvent extraction of the product chiral amine, comprising a so-called reaction phase and a so-called stripping phase at different pH, separated by a supported liquid membrane, in casu a membrane with the pores loaded/impregnated with undecane. In the reaction phase (at pH≥9), the transaminase reaction takes place with isopropyl amine (IPA) or alanine (ALA) as amine donor. Although all uncharged substrate and product compounds can be (selectively) extracted by the organic solvent of the liquid membrane and can pass through the liquid membrane to the stripping phase (at pH<3), back extraction from the stripping phase into the organic membrane phase is prevented for the charged compounds (in casu the non-aminoacid amine donor and amine product in the stripping phase). This setup results in a selective in situ product removal, which steers the reaction equilibrium towards chiral amine synthesis. However, the liquid-membrane setup used in this study is not stable due to the liquid in the membrane pores being prone to leaking out during operation, requiring regularly regeneration of the liquid membrane, and is sensitive to transmembrane pressure differences. This setup is thus less suitable for long-term, large scale chiral amine synthesis and separation.\nWO01/09042 discloses a process for recovering an aromatic amine dissolved in an aqueous fluid comprising transferring the undissociated aromatic amine from the aqueous fluid to an acidic stripping solution across a membrane, wherein the membrane is a non-porous, elastomeric (polymeric) selectively permeable membrane. In the process of WO01/09042, the uncharged/undissociated aromatic amines pass through the membrane by diffusion to the acidic stripping phase, which can be considered as the dissolution of the undissociated aromatic amines within the polymeric membrane. Similar as in Börner et al (Organic Process Research & Development, 2015, 19(7), 793-799), back extraction from the acidic stripping phase is prevented for the charged compounds. However, the performance and stability of an elastomeric polymeric membrane of WO01/09042 is often negatively affected when using particular solutions or liquids in the separation process or during membrane cleaning.\nThere remains a need in the art for the improved separation of organic compounds such as amines or organic acids from a solution, particularly from a solution comprising a chiral amine and/or multiple amines, that overcomes the limitations of the prior art. In this context, there is a particular need for improved enzymatic processes and systems to produce chiral amines that overcome at least some of the above indicated limitations of the transamination reaction and at the same time allow for long term application, which can be adapted to specific substrate donor amines/ketones and/or product amines."}
{"text": "1. Field of the Invention\nThe present invention relates to a light emission device, a method of manufacturing the light emission device, and a display device using the light emission device as a backlight unit. More particularly, the present invention relates to a driving electrode of the light emission device.\n2. Discussion of Related Art\nA liquid crystal display, which is one of a variety of flat panel display devices, displays an image by varying a light transmission amount at each pixel using a dielectric anisotropy property of liquid crystal whose twisting angle varies according to a voltage applied.\nThe liquid crystal display includes a liquid crystal panel assembly and a backlight unit for emitting light toward the liquid crystal panel assembly. The liquid crystal panel assembly receives light emitted from the backlight unit and allows the light to be transmitted or blocked by a liquid crystal layer.\nBacklight units can be classified according to their light sources into different types (or light source types), one of which is a cold cathode fluorescent lamp (CCFL) type. The CCFL is a linear light source that can uniformly emit light to the liquid crystal panel assembly through optical members such as a diffusion sheet, a diffuser plate, and a prism sheet.\nHowever, in the CCFL type backlight unit, since the light emitted from the CCFL travels through the optical members, there may be a light loss. Considering the light loss, a relatively high intensity of light is emitted from the CCFL. This causes an increase in power consumption. Furthermore, since it is difficult to increase the size of the CCFL type backlight unit due to a structural limitation of the CCFL type backlight unit, the CCFL type backlight unit may not be applied to a large-size liquid crystal display (e.g., a liquid crystal display that is over 30 inches).\nA light emission diode (LED) type backlight unit is another light source type that may be used to provide light to a liquid crystal display. The LED type backlight unit includes a plurality of LEDs and optical members such as a reflection sheet, a waveguide plate, a diffusion sheet, a diffuser plate, a prism sheet, and the like. The LED type backlight unit has fast response time and relatively good color reproduction characteristics. However, the LED is costly and increases an overall thickness of the liquid crystal display.\nTherefore, recently, a field emission type backlight unit that emits light using an electron emission property of an electric field has been developed to replace the CCFL and LED type backlight units.\nSince the field emission type backlight unit should provide a relatively high luminance, a relatively high voltage is applied to an anode electrode of the field emission type backlight unit. As such, there may be a short circuit between the driving electrode formed on a first substrate of the field emission type backlight unit and the anode electrode formed on a second substrate of the field emission type backlight unit. In addition, an arc discharge may be generated in a vacuum envelope (or chamber) of the field emission type backlight unit. The arc discharge may damage the driving electrodes.\nFurthermore, among the driving electrodes, cathode and gate electrodes should be insulated from each other. However, the cathode and gate electrodes may be short-circuited during a manufacturing process.\nWhen the driving electrodes are short-circuited or damaged as described above, the pixels controlled by the driving electrodes cannot properly emit light. That is, in a conventional field emission type light emission device (or backlight unit), when one driving electrode is damaged, all of the pixels that are controlled by the damaged driving electrode are defective in the emission of light.\nIn addition, the conventional backlight units are maintained at a constant (or predetermined) brightness when the liquid crystal display is driven. Therefore, it is difficult to improve the display quality of the liquid crystal display to a sufficient level.\nFor example, when the liquid crystal panel assembly intends to display an image having a high luminance portion and a low luminance portion in response to an image signal, it will be possible to realize an image having a more improved dynamic contrast if the backlight unit can emit lights having different intensities to the respective high and low luminance portions.\nHowever, since the conventional backlight units may not achieve the above function, the liquid crystal display has a limitation in how much the dynamic contrast of an image can be improved."}
{"text": "1. Field of the Invention\nThe present invention relates to a controller and a machining apparatus having the same.\n2. Description of the Related Art\nThere has been known a technique of detecting the position of an object based on two periodic signals having a phase difference of 90° which are generated by an encoder for encoding an object position. A count value obtained by counting the periods of one of the two periodic signals indicates an object position in periods. On the other hand, the arctangent value of the ratio of the two periodic signal values indicates an object position during one period. A signal representing the final object position can be obtained by combining the count value and the arctangent value. Unless the period count operation and the arctangent value processing operation completely synchronize, periodic errors occur in the value obtained upon combining the count value and the arctangent value. Conventionally, processing of compensating the periodic errors is done based on only the current values of the two periodic signals.\nHowever, as the movement of the object speeds up with respect to the sampling period, it becomes more difficult to compensate the above-described periodic errors caused by the lack of synchronism between the period count operation and the arctangent value processing operation."}
{"text": "This invention relates to a waterwall and, more particularly, to a waterwall that has been provided with a surface that is resistant to erosion caused by particulate material.\nMost erosion in industrial processes occurs where large volumes of abrasive particles such as coal, catalyst, sand, shale, limestone, etc. change direction of flow via elbows, tees, separators, swirl vanes and the like. The erosion occurs as a result of low angle impingement by large volumes of the abrasive particules which move at varying velocities.\nSeveral systems have been devised to resist this type of erosion. For example, monolithic, cement or phosphate bonded castable (and plastic) refractories held by steel anchors on about two or four-inch staggered centers, or on a hexagonal steel grid have been utilized to resist the above erosion. Also, a considerable thickness, up to several inches of the above mentioned refractories, have been installed on V-bar or S-bar anchors, and if the erosion is exceptionally severe a prefired refractory is used which is usually bolted to the steel structure.\nHowever, the use of these refractories increase the thermal conductivity which, in many systems, seriously affects heat absorption rates to the tubular surfaces in fluidized bed boilers, for example.\nIt has been discovered that a steel support system is more erosion resistant than most of the monolithic cement-bonded, castable refractories discussed above that are traditionally used in studded anchor wall systems. However, if the tube enclosure walls were simply designed with a greater thickness of steel, the weight and cost increase could be prohibitive."}
{"text": "1. Field of the Invention\nThe present invention relates to a singlet oxygen generating process and also to a process of oxidizing fullerenes with the singlet oxygen.\n2. Description of Related Art\nOxidation of fullerenes is a key for synthesis of organic compounds. Oxidation of fullerenes with molecular oxygen is conventionally performed by a photo-chemical process using organic dyes as photo-sensitizers. Use of the molecular oxygen as an oxidant has been regarded as promising means for clean reactions with less environmental burdens, compared with use of toxic organic or inorganic oxidants.\nFullerenes, which are recently applied to various technical fields, are also effective photo-sensitizers, as reported in Ref. No. 1. According to Ref. No. 1, singlet photo-excited fullerenes are further excited to a triplet state by intersystem crossing. The excitation is accompanied with energy transfer to oxygen, and oxygen is sequentially excited to a singlet state.\nBy the way, fullerenes represented by C60, C70 or C84 are carbon clusters with closed shell structures (hereinafter, referred to as “carbon cages”) and have peculiar physical properties originated in the specified carbon cages. For instance in a medical field, big powers for generation of singlet oxygen enable employment of fullerenes as anti-cancer drugs or anti-HIV proteases in combination with selective drug-delivery systems. In a cosmetic field, fullerenes are used as anti-active oxygen medicines similar to vitamin C due to radical-repairing powers. Moreover, development of fullerenes is also researched and examined by bonding other elements or groups to some of cage-forming carbon atoms or by enclosing other elements or groups in the carbon cages so as to change properties of fullerenes themselves.\nAmong oxidation processes using fullerenes as photo-sensitizers, oxidation of fullerenes is the simplest system. Since fullerene oxide has excellent activity, its applicability as an additive to battery electrodes using LUMO (a lowest unoccupied molecular orbital) lower than fullerenes is also expected, as well as precursors for functional fullerene derivatives. For instance, fullerene oxides are produced by liquid-state photo-oxidation as noted in Ref. No. 2 or by chemical oxidation with oxidants, e.g. m-chloroperbenzoic acid, as noted in Ref. No. 3. Ref. No. 4 reports another fullerene oxidizing process, wherein fullerenes are irradiated with microwaves in a solid state mixed with m-chloroperbenzoic acid.                Ref. No. 1 Michael Orfanopoulous, Spiros Kambourakis, Tetrahedron Lett. (1994) 35, 1945        Ref. No. 2 Kathleen M Creegan, John L. Robbins, Win K. Robbins, John M. Millar, Rexford D. Sherwood, Paul J. Tindall, Donald M. Cox, J. Am. Chem. Soc. (1992) 114, 1103        Ref. No. 3 Alan L. Balch, David A. Costa, Bruce C. Noll, Marilyn M. Olmstead, J. Am. Chem. Soc. (1995) 177, 8926        Ref. No. 4 Weon Bae Ko, Sung Ho Hwang, Ju Hyun Ahn, Elastomer, (2005) 40, 45.        \nAny method of Ref. Nos. 1-3 premises liquid-state reactions, which are unavoidably accompanied with heavy environmental burdens and unsuitable for massive synthesis of fullerene oxides since huge amounts of solvents are consumed for dissolution of sparingly-soluble fullerenes. Besides, photo-irradiation is essential for the reactions.\nThe other method of Ref. No. 4 is based on solid-state reactions and limits a reaction field between solids. As a result, objective compounds are produced with poor yields, and organic or inorganic oxidants were still required."}
{"text": "Cache partitioning to serve multiple client processes (“clients”) in computing systems has several benefits. Partition size can vary to match performance requirements of the client. Clients experience robust performance since data in its private cache is not subject to evictions due to other clients' accesses. Idle partitions can be turned off to reduce power consumption. However, these benefits are not as great as they could be due to limitations imposed by the way current cache partitioning is accomplished.\nCurrently caches may be partitioned across cache ways only, which has a number of limitations. The cache partition number is limited by the number of cache ways, and adding associativity to the cache partitions to reduce cache conflicts reduces the maximum number of cache partitions possible. Minimum way partition size is determined by the associativity. For example, the smallest partition for an 8 MB cache with 8 ways that is 8-way associative is 1 MB. Also, smaller cache partitions could result in loss of cache space due to fragmentation. Using the same example, 512 KB of data stored in the smallest 1 MB cache partition would result in 512 KB of unused space in the cache partition. Minimum way partition power collapse granularity is determined by the associativity. In the same example, the minimum granularity that can be power collapsed is a 1 MB cache partition. The size of collapse granularity affects power savings/usage, especially if a shared system cache replaces multiple private memories that can be independently power collapsed.\nAnother problem faced in mobile system on chip (SoC) computing systems is the diverse access patterns generated by subsystems. For example, the multimedia subsystems, like video and graphics, have longer access times and benefit from longer cache lines, while the accesses from the applications central processing unit (CPU) are more random and sensitive to latency and thus benefit from smaller cache lines and higher associativity. Not all of the potential access patterns work well with current cache partitioning that relies on way cache partitioning. Conventional cache partitioning provides one-size caches for all applications and subsystems. The size of the cache partitions may benefit certain access patterns more than others. Another example is the producer-to-consumer access pattern seen between an image processing engine and a display engine, which may benefit from a first in first out (FIFO) replacement policy. However, a cache holding the transit data between the engines may use an inefficient replacement policy for the application."}
{"text": "In modern \"enterprise\" digital data processing systems, that is, computer systems for use in an office environment in a company, a number of personal computers, workstations, and other devices such as mass storage subsystems, network printers and interfaces to the public telephony system, are typically interconnected in a computer network. The personal computers and workstations (generally, \"PCs\") are used by individual users to perform processing in connection with data and programs that may be stored in the network mass storage subsystems. In such an arrangement, the PCs, operating as clients, access the data and programs from the network mass storage subsystems for processing. In addition, the PCs will enable processed data to be uploaded to the network mass storage subsystems for storage, to a network printer for printing, to the telephony interface for transmission over the public telephony system, or the like. In such an arrangement, the network mass storage subsystems, network printers and telephony interface operate as servers, since they are available to service requests from all of the clients in the network. By organizing the network in such a manner, the servers are readily available for use by all of the PCs the network. Such a network may be spread over a fairly wide area, with the PCs being interconnected by communication links such as electrical wires or optic fibers.\nA problem arises when an operator at one PC wishes to determine what other PCs are connected to and available for use over the network, or to determine what other resources are available over the network. For example, an operator of one PC connected to the network may wish to use programs or data which are on other PCs (including both client PCs and server PCs) which are connected to the network, in which case he or she may need to initially obtain the identifications of the particular PCs which are connected to the network. After obtaining the identifications, the operator may retrieve information from or transfer information to their respective information storage systems. The operation by which the operator initially identifies the personal computers, workstations and other resources which are available over the network is typically referred to as \"browsing.\"\nSeveral general methodologies have been used to enable an operator to determine the various PCs and other resources which are available over the network. In one methodology, each PC can be provided with one or more lists of other PCs and other resources which are available to it over the network. In such a case, the list would need to be updated by, for example, a system administrator as each PC or other resource is added to or removed from the network. There are several problems with this methodology. First, maintaining all of the PC/resource lists in an updated condition can be a difficult task, particularly as the size of the network increases and PCs and other resources are added or removed in remote portions of the network. In addition, such lists would normally only identify the particular PCs and resources which are connected to the network, and would not necessarily indicate whether they are available for use at any particular point in time; thus, for example, if a particular PC or other resource is switched off, it would still be identified in a list but would not actually then be available.\nIn another methodology, each PC or other resource periodically broadcasts \"advertising\" messages containing their respective identifications over the network, perhaps along with other information such as various types of resources that they may provide which may be accessed over the network. The various PCs can receive the broadcast messages and from the information contained therein can identify the PCs and other resources which broadcast the respective messages. The PCs can buffer the information that they receive in the respective broadcast messages and provide it to the operator when requested. Since new PCs may be added to the network, the PCs and other resources which broadcast the messages will need to broadcast them periodically over the network. In addition, if the PCs receiving the messages determine that they have not received broadcast messages from a particular PC or other resource within a particular time interval, they (that is, the receiving PCs) can determine that the particular PC or other resource has, for example, been turned off or is otherwise no longer available. While this methodology allows each PC to identity other PCs and resources which are available over the network, and maintain the information relatively current, without requiring assistance from a system administrator or other operator, it does require the PCs and other resources to broadcast a number of messages. A number of the broadcast messages may be unnecessary, either because the PCs and other resources which are connected to and available over the network has not changed, or because operators have not requested information as to the PCs and resources which are available over the network. This problem is exacerbated as the number of PCs and other resources which are connected to the network increases, since the number of advertising messages would increase correspondingly and reduce the network bandwidth which is available for transfer of other information over the network. Additionally for networks which are interconnected by devices such as routers, the broadcast advertising messages may not be transferred throughout the entire network, since typically routers are designed so as not to forward broadcast messages.\nA further methodology has been developed which alleviates these problems to some extent. In this further methodology, when an operator at a PC wishes to know what other PCs and other resources are available over the network, his or her PC will broadcast an inquiry message over the network. Other PCs and resources connected to the network are to respond to the inquiry message with the required browsing information, including, for example, their respective identifiers, within a particular time-out period. The PC that had broadcast the inquiry message will receive the browsing information can provide the information which it receives during the time-out period to the operator. Generally, to limit the information provided, the broadcast message is responded to only by PCs and other resources within a particular \"workgroup\" to which the requesting PC has been assigned, so that the information displayed will be only the information for the particular workgroup. While this methodology alleviates the transmission of unnecessary messages which will be broadcast in the previously-described methodology, and it ensures that the information provided to the operator is current, it does have several problems. First, the assignment of PCs and other resources to a particular workgroup may be unduly limiting. In addition, the requirement that a PC or other resource respond within the time-out period in order for its information to be included in the information displayed to the operator, results in information not being included if it is not received within the required time-out period. Thus, if a PC or other resource is busy and not able to respond to the broadcast message within the time-out period, its information will not be provided to the operator even though it may be available for use."}
{"text": "1. Field of the Invention\nThe present invention relates generally to dolls and miniature buildings and, more particularly, to a doll and a miniature building that simulates a real-life person and building.\n2. Description of the Prior Art\nDolls and doll houses or other buildings scaled down to doll size are well known in the art. Aside from doll houses, it is well known to have miniature farms, miniature battlefields, and miniature gasoline stations, among others, to be associated with corresponding dolls. An environment that the prior art of dolls and miniature buildings does not disclose is a ballet studio. Many persons, both children and adults, spend a considerable amount of time in a ballet studio. In this respect, for a number of reasons, it would be desirable if a doll and a simulated building were provided which represented a ballerina and a ballet studio. Some reasons for the desirability of a simulated ballerina and dance studio are as follows. They would serve as a reminder to a ballet student to practice ballet exercises. They could serve as an attractive decoration in a ballet studio.\nWhen a ballet student leaves the ballet studio, it is often difficult for the student to remember the body positions that are required. In this respect, it would be desirable if a simulated ballerina were provided that could be adjusted to depict body positions of a ballerina which may be difficult for the student to remember.\nIn a ballet class, a number of exercises are done in association with a horizontal bar in front of a wall mirror. The horizontal bar is used for balancing the student as the student goes through a number of exercises. In this respect, it would be deskable if a simulated ballerina and dance studio apparatus were provided which included a horizontal bar onto which the hand of the ballerina doll can be placed in order to balance the doll.\nThroughout the years, a number of innovations have been developed relating to doll houses, and the following U.S. patents are representative of some of those innovations: U.S. Pat. Nos. 3,888,039; 4,030,234; 4,070,789; 4,883,443; and 4,919,982. More specifically, the above-cited patents relate to doll buildings that can be folded or disassembled. However, none of the patents disclose a simulated ballerina or a simulated ballet studio.\nMany dolls are known in the prior art. Some of the prior art dolls include articulated joints. However, none of the dolls in the known prior art contain the number and quality of articulated joints that would permit a ballerina doll to be adjusted into known ballet positions. In this respect, it would be desirable if a doll were provided which contained the number and quality of articulated joints that permit the ballerina doll to be adjusted into known ballet positions.\nAn articulated joint that permits the ballerina doll to be adjusted into known ballet positions should be able to be easily adjusted and should retain its adjusted position.\nIn a simulated ballet studio there are a number of features that would be most significant in the simulation. The key features to be simulated are the horizontal bar, the wall mirror, and the dance floor.\nThus, while the foregoing body of prior art indicates it to be well known to use dolls and doll houses, the prior art described above does not teach or suggest a simulated ballerina and dance studio apparatus which has the following combination of desirable features: (1) represents a ballerina and a ballet studio; (2) can be adjusted to depict body positions of a ballerina which may be difficult for a ballet student to remember; (3) includes a horizontal bar onto which the hand of a ballerina doll can be placed in order to balance the doll; (4) contains the number and quality of articulated joints that permit the ballerina doll to be adjusted into known ballet positions; (5) includes articulated joints that permit the ballerina doll to be adjusted into known ballet positions easily and with good retention of the adjusted position; and (6) has a simulated ballet studio which includes a horizontal bar, a wall mirror, and a dance floor. The foregoing desired characteristics are provided by the unique simulated ballerina and dance studio apparatus of the present invention as will be made apparent from the following description thereof. Other advantages of the present invention over the prior art also will be rendered evident."}
{"text": "The invention relates to glycoprotein B segments of herpes simplex virus types 1 and 2 (HSV-1, HSV-2) containing linear antigenic epitopes reactive with human antibodies to HSV-1 and HSV-2. Of particular interest are HSV glycoprotein B segments containing type-specific epitopes useful in serodiagnostic immunoassays for distinguishing HSV-1 infection from HSV-2 in humans, and in human vaccines for generating neutralizing antibodies to HSV-1 or HSV-2.\n1. Field of Art\nThe herpes simplex virus (HSV) glycoprotein B (gB) is a transmembrane envelope glycoprotein that contributes to the penetration of virions into host cells by host cell recognition, viral adherence, and virion-cell membrane fusion. gB induces fusion of the virion envelope with the cellular cytoplasmic membrane, an essential function for virus entry. Immunization of humans and animals with purified gB induces virus-neutralizing antibody responses, and recombinant gB proteins are being investigated as HSV vaccines. The amino acid sequences for both gB-1 and gB-2 and their encoding nucleotide sequences are known and available, e.g., from Genebank.\nThe HSV type 2 gB (gB-2) polypeptide contains 904 amino acids (aa)--(FIG. 5). The amino-terminal segment (1-22aa) is a signal peptide that is cleaved from the mature protein during processing. aa 745 to 798 constitute a hydrophobic transmembrane anchor domain. The segment carboxy-proximal to the anchor domain (aa 799 to 904) is cytoplasmic and appears to mediate membrane fusion. The amino-proximal segment (aa 23 to 744) is extracellular and contains epitopes that are recognized by virus-neutralizing antibodies. HSV type 1 gB (gB-1) is structurally similar (FIG. 4).\nHSV-2 and HSV-1 are closely related viruses. Most of their proteins, including gB, are highly conserved (very homologous) and are known to elicit cross-reactive antibody responses. It has accordingly been difficult to provide reliable, sensitive immunoassays capable of distinguishing HSV-1 from HSV-2 infection, or to provide type-specific vaccines effective against HSV infection.\n2. Discussion of Related Art\nHSV infections in humans are commonly diagnosed by immunoassay of blood samples for HSV antibodies in Western blot assays or other immunoassays using random lysates of cells infected with HSV-1 or HSV-2 as antigen targets. These serodiagnostic assays have been recognized as generally unsatisfactory for distinguishing HSV-1 infection from HSV-2 infection owing to strong responses of type-common antibody-reactive regions of both native proteins. This has proved to be a particular problem in distinguishing acute HSV-2 infection in patients previously infected with HSV-1, because of a likely anamnestic response to HSV-1 and HSV-2 type-common antigens which obscures responses to type-specific antigens. Difficulties in detecting HSV-2 type-specific antibodies by using HSV-2 cell lysate Western blot assays on subjects previously infected with HSV-1, for example, may result in a not uncommon misclassification of samples subjected to these assays as HSV-1 positive and HSV-2 negative. Accordingly, on-going research has attempted to identify domains of various native HSV viral proteins which are highly reactive with anti HSV-1 and anti HSV-2 antibodies and which provide a virus type-specific antibody response.\nWhile several HSV proteins have been considered as sources of useful antigenic epitopes for diagnosis and prevention, only HSV-1 and HSV-2 glycoprotein G (gG-1, gG-2) has been found useful to date for use in assays with improved sensitivities and specificities. Most HSV-2 and HSV-1 proteins are very homologous, as noted above. However, HSV-2 gG-2 and HSV-1 gG-1 have highly dissimilar amino acid sequences, and human gG-2 and gG-1 antibody responses in the presence of HSV-2 and HSV-1 infections appear to have significantly improved virus type-specificity over lysates of infected cells. Thus, detection of antibody reactivities to native or recombinant HSV-1 or HSV-2 glycoprotein G presently forms the basis for current HSV type-specific immunoassays.\nHSV glycoprotein B has also sparked interest as a possible source of type-specific HSV antigens. To date, however, only the complete proteins have been used, and these have exhibited only cross-reactive antibody responses. Virus-specific determinates to human antibodies have been eagerly sought for many years, without notable success.\nHSV-2 and HSV-1 infections each elicit strong immunoglobulin G (IgG) antibody responses in vivo to both HSV-2 gB-2 and to HSV-1 gB-1 (1, 2, 3); these antibodies are detectable early in acute infections, attain high titers, and persist for many years (4, 2, 5). Cross-reactive polyclonal antibody responses predominate (4, 6, 7). This exhibition of the related antigenicity of HSV gB-2 and gB-1 is consistent with the marked overall conservation of their amino acid sequences (86% at the amino acid level). Perhaps in consequence, the structures of gB-1 and gB-2 have not been previously broadly elucidated with respect to specific antigenic domains recognizing human antibodies which are clinically useful, especially type-specific antigenic regions capable of distinguishing HSV-1 from HSV-2 infection.\nWhile human IgG antibodies cross-reactive to whole or randomly lysed gB-2 and gB-1 are thus known, virus type-specific antibodies have only been reported in murine models. Immunization of laboratory animals with purified virion-derived gB-2 elicits antibodies that neutralize HSV-2 infectivity in vitro and protect susceptible animal hosts from experimental HSV-2 infections (reviewed in Seminars in Pediatric Infections 2: 178-185, Stanberry, 1991 and Microbiol. and Immunol. 179: 137-158, Burke, 1992).\nHSV-2 infections and HSV-1 infections elicit strong IgG antibody responses to gB-2 and to gB-1, respectively. gB-2 responses include antibodies that cross-react strongly with native gB-1, and gB-1 responses include antibodies that cross-react strongly with native gB-2. Therefore, gB-2 reactivities and gB-1 reactivities do not differentiate between antibody responses to HSV-2 infections and antibody responses to HSV-1 infections. Immune responses to most HSV-2 proteins include antibodies that cross-react with the homologous HSV-1 protein, and vice versa. Linear epitopes that react with type-specific, virus-neutralizing murine monoclonal antibodies have been localized to the extracellular amino-terminal segment of gB-1 (8, 9, 10, 11, 12, 13, 14). However, little information has been available regarding the locations and type specificities of gB epitopes recognized by human HSV antibodies (8, 15, 14). Vaccines comprising both recombinant gB-2 and gB-2 are currently undergoing clinical trials (16, 17), but efficacy or type-specificity of gB-2 in this application is not yet known.\nThe above numerals refer to the following publications, all incorporated herein by reference:\n1. J. Virol. Methods 18:159-168, 1987 PA0 2. Infect. Immun. 31:1062-1070, 1981 PA0 3. Am. J. Epi 104:192-201, 1976 PA0 4. J. Med. Virol. 17:153-166, 1985 PA0 5. Infect. Immunol. 34:880-887, 1981 PA0 6. J. Med. Virol. 15:251-263, 1985 PA0 7. Rev. Infect. Dis. 2:899-913, 1980 PA0 8. J. Med. Virol. 27:309-316, 1989 PA0 9. J. Gen. Virol. 70:735-741, 1989 PA0 10. Virol. 135:379-394, 1984 PA0 11. Virol. 186:99-112, 1992 PA0 12. Virol. 172:11-24, 1989 PA0 13. J. Infect. Dis. 166:623-627, 1992 PA0 14. J. Infect Dis. 168:844-853, 1993 PA0 15. J. Clin. Micro. 23:725-730, 1986 PA0 16. Rev. Infect. Dis. 13:S906-S911, 1991 PA0 17. Micro. Immunol. 179:137-158, 1992"}
{"text": "1. Field of the Invention\nThe present invention relates to a signal conversion method and apparatus, a recording medium, and a print system, and more particularly, to a signal conversion technique for limiting the amount of color material used in a print system that forms an image on a recording medium using color materials.\n2. Description of the Related Art\nA printing apparatus (printing device), such as an ink jet printer or an electrophotographic printer, applies a color material (recording agent), such as ink or toner, onto a recording medium, such as a sheet, to form a desired image on the recording medium.\nJP2005-101934A discloses a technique which performs an under-color removal process of replacing a predetermined amount of CMY signal, which is an input CMYK signal, with a K signal in order to prevent excess color material from being applied and to improve the reproducibility of an output image while limiting the amount of toner or ink attached to a unit area (paragraph 0005 in JP2005-101934A).\nThat is, JP2005-101934A discloses “an image processing apparatus that supplies an image signal to an image formation apparatus for forming an image on a recording medium using a recording agent and includes: a storage unit that stores the limit value of the total sum of the levels of image signals output to the image formation apparatus; input means for inputting CMYK signals as the image signals; first comparison means for comparing the total sum of the levels of the input CMYK signals with the limit value; under-color removal means for performing an under-color removal process of replacing a predetermined amount of CMY signals with a K signal when the total sum of the levels of the input CMYK signals is greater than the limit value; second comparison means for comparing the total sum of the levels of the CMYK signals subjected to the under-color removal process with the limit value; and an adjustment unit that reduces CMY signals after the under-color removal process by the same proportions such that the total sum after the under-color removal process is equal to or less than the limit value when the total sum after the under-color removal process is greater than the limit value” (claim 1 of JP2005-101934A)."}
{"text": "This application claims the benefit of priority from prior Japanese Patent Application P2001-303086 filed on Sep. 28, 2001; the contents of which are incorporated by reference herein.\nI. Field of the Invention\nThe present invention relates to a semiconductor device having an oscillating circuit. More particularly, the invention relates to a semiconductor device having a frequency oscillating circuit, such as an LC (which means xe2x80x9cinductive-capacitivexe2x80x9d) oscillating circuit.\nII. Background of the Invention\nThe following will explain an example of the arrangement of a conventional oscillating circuit. FIG. 1 is a circuit diagram illustrating the conventional oscillation circuit arrangement.\nIn the oscillating circuit of FIG. 1, transistors Q101 and Q102 form a differential pair, and each base of the transistors Q101 and Q102 is connected to a collector of the other transistor. A coil L101 and a capacitor C101, which form an LC resonant tank circuit, are connected between the collectors of these transistors Q101 and Q102.\nThe collectors of transistors Q101 and Q102 are connected to a power supply voltage VDD via resistors R101 and R102. A constant current source Io is connected to the emitters of transistors Q101 and Q102. Moreover, the constant current source Io is connected to the ground potential GND. The output signal OUT1 from the collector of the transistor Q101 and an output signal OUT2 from the collector of the transistor Q102 are used as a differential pair output, or, alternatively, a single output of one side of the oscillating circuit (e.g., OUT1) is used.\nHowever, the aforementioned oscillating circuit can produce a signal at a fixed frequency, and has a problem in which the frequency cannot be varied.\nAccordingly, the present invention is directed to a semiconductor device having a frequency oscillating circuit, such as an LC oscillating circuit.\nA semiconductor device according to an embodiment of the present invention includes, for example, a first transistor having a control electrode and having one end and the other end of a current path thereof; a second transistor having a control electrode and having one end and the other end of a current path thereof, the control electrode is coupled to one end of the current path of the first transistor and one end of the current path thereof is coupled to the control electrode of the first transistor; a current mirror circuit that supplies a current to one end of the current path of the first transistor and one end of the current path of the second transistor; an inductor coupled to one end of the current path of the first transistor and one end of the current path of the second transistor; a first capacitor coupled to one end of the current path of the first transistor and one end of the current path of the second transistor; and a second capacitor and a switch element, are coupled to one end of the current path of the first transistor and one end of the current path of the second transistor.\nA semiconductor device according to another embodiment of the present invention includes, for example, a first transistor having a control electrode and having one end and the other end of a current path thereof; a second transistor in which a control electrode is coupled to one end of the current path of the first transistor and one end of a current path thereof is coupled to the control electrode of the first transistor; a current mirror circuit that supplies a current to one end of the current path of the first transistor and one end of the current path of the second transistor; an inductor coupled to one end of the current path of the first transistor and one end of the current path of the second transistor; a constant current source coupled to the other end of the current path of the first transistor and the other end of the current path of the second transistor; a first capacitor coupled to one end of the current path of the first transistor and the other end of the constant current source; a second capacitor coupled to one end of the current path of the second transistor and the other end of the constant current source; a third capacitor and a first switch element coupled to one end of the current path of the second transistor and the other end of the constant current source; and a fourth capacitor and a second switch element coupled to one end of the current path of the second transistor and the other end of the constant current source.\nOther features and advantages of the present invention will become apparent from the following detailed description. It should be understood, however, that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the spirit and scope of the invention will become apparent to those skilled in the art from this detailed description."}
{"text": "U.S. Pat. Nos. 6,386,295 and 6,942,430, which are incorporated here by reference, disclosed the use of vibratory and rotary drivers for the installation of pipe piling. Pipe piles, as used in the installation of structural foundations or geothermal piles, are segments of pipe that must be connected and driven together from the surface to reach the desired depth. Consequently, whether used in connection with vibratory or rotary drivers, the connection between pipe pile segments is vitally important to maximizing the driving power and reducing wear and tear in the pipe segment connection points. As the length of the column increases, weaknesses in the junctions between the pipe pile segments weaken the entire column, making it important to limit movement in the junctions.\nThus prior art methods that require the use of bolts through pipe pile and connectors increase the wear on the piles and the fasteners as some of the power from the drivers is transferred inefficiently through such joints. U.S. Pat. Nos. 6,386,295 and 6,942,430 disclosed pipe pile coupling mechanisms that may render the use of fasteners unnecessary, however, the force applied to the coupled joint through threading also created wear in the coupler and the pipe pile due to over threading or thread damage.\nWhile it is generally acknowledged that installation of pipe piling is improved in stability and/or strength when installed with grout or similar material along the exterior of the column, prior art methods, including those disclosing push-out tips, are limited by the actual ability to push out the tip at the bottom of the column, or by difficulty in handling the grout during installation.\nThe citation of documents herein is not to be construed as reflecting an admission that any is relevant prior art. Moreover, their citation is not an indication of a search for relevant disclosures. All statements regarding the date(s) or contents of the documents is based on available information and is not an admission as to their accuracy or correctness."}
{"text": "High molecular weight linear polyesters and copolymers of glycols and terephthalic or isophthalic acid have been available for a number of years. These are described inter alia in Whinfield, et al, U.S. Pat. No. 3,047,539. These patents disclose that the polyesters are particularly advantageous as film and fiber-formers.\nPoly(1,4-butylene terephthalate), because of its very rapid crystallization from melt, is uniquely useful as a component in injection moldable compositions. Work pieces molded from such polyester resins, in comparison with other thermoplastics, offer a high degree of surface hardness and abrasion resistance, high gloss and low surface friction.\nOne useful family of such compositions comprises those which are reinforced, e.g., with from about 5% to about 55% filamentous glass, based on the weight of glass and polyester components, and optionally rendered flame retardant by the incorporation of a flame-retarding amount of a flame retarding compound.\nThe versatility of injection moldable compositions containing poly(1,4-butylene terephthalate) resin has been demonstrated, however, it is always desirable to increase the versatility of these compositions by the addition of additives and the like to modify the physical properties thereof. For example, U.S. Pat. No. 4,104,242, Kochanowski, et al., disclose the addition of a hydrazine compound to poly(1,4-butylene terephthalate) resin for modifying resistance to high voltage breakdown. Other properties, including impact strength, flexural strength, flexural modulus, tensile strength, melt viscosity, arc track rate, etc., may also be the object of such improvement efforts.\nIt is an object of the present invention to improve the physical properties of articles molded from glass fiber reinforced poly(1,4-butylene terephthalate) resin molding compositions.\nIt is another object of the present invention to improve the physical properties of articles molded from fiber reinforced, flame retarded poly(1,4-butylene terephthalate) resin molding compositions."}
{"text": "Considerations of energy saving and reduction of pollution, particularly in an urban environment, are leading motor vehicle manufacturers to equip their models with an automatic system for stopping/re-starting the thermal engine (known as the “automatic engine stop/re-start system” hereinafter in the description), such as the system known by the term “stop and go”.\nAs recalled by the company VALEO EQUIPEMENTS ELECTRIQUES MOTEUR in patent application FR2875549, it is possible for vehicles to function according to the “stop and go” mode by means of a reversible electrical machine, or alternator-starter, which is coupled to the thermal engine, and is supplied by an inverter in “starter” mode.\nThe use of an alternator-starter system in a “stop and go” functioning mode consists in certain conditions of giving rise to complete stoppage of the thermal engine when the vehicle itself is at a standstill, then of re-starting the thermal engine, as the result for example of action by the driver which is interpreted as a request to re-start the thermal engine.\nA typical “stop and go” situation is that of stoppage at a red light. When the vehicle stops at the light, the thermal engine is automatically stopped, then, when the light turns to green, the engine is re-started by means of the alternator-starter system, as the result of diagnosis, by the system, of the clutch pedal being pressed down by the driver, or any other action which represents the wish of the driver to give rise to movement of his vehicle.\nFor this purpose, as the company VALEO EQUIPEMENTS ELECTRIQUES MOTEUR also indicates in patent application FR2875551, the “stop and go” system uses information representative of the state of functioning of the vehicle which is obtained from sensors of the vehicle, such as temperature sensors in the passenger space, or detection of the position of the clutch pedal, or information read on a data communication bus, such as the speed of the vehicle or the speed of rotation of the engine.\nIt will be appreciated that the automatic re-starting function carried out by an alternator-starter system is a function which must assure a certain level of safety of functioning.\nAs described by the company VALEO EQUIPEMENTS ELECTRIQUES MOTEUR in patent application FR287557, the essential characteristic of a method for controlling an alternator-starter system in starter mode is that it must comprise steps of verification of so-called prior conditions for re-starting of the engine before the re-starting phase takes place, and so-called “safety” conditions throughout the phase of re-starting of the engine.\nA prior condition for restarting of the engine is determined in particular by at least one intentional action of the driver, for example pressing the clutch pedal down by more than 90%.\nAn additional condition is, for example, that the network on board the vehicle must be supplied, i.e. the contact key of the vehicle must be in the position +APC “after contact”, and the traction chain must be open.\nIn fact, when the vehicle is at a standstill (i.e. with zero speed) and the thermal engine is stopped by the automatic function, there is no question of re-starting the engine unless it is certain that the traction chain is open.\nIf the chain were closed, re-starting the engine could impel the vehicle, and consequently make it advance (or reverse), thus endangering the safety of property or people.\nThus, in order to avoid this type of situation, a solution is now implemented which uses a specific control line Auth. Ond., in order to authorise the control of the inverter. This specific control line Auth. Ond. is derived from a control system of the vehicle, which is external to the automatic stop/re-start function, thus fulfilling most of the safety requirements.\nIn addition, in order to detect the activated state of the electrical machine, a supplementary safety device is introduced. This safety device is constituted by the line Diag.1 Ond., which conveys information which indicates the rotation of the electrical machine, and is obtained from signals from position sensors of the rotor.\nThis information is then put into form by the microprocessor for control of the electrical machine, and is transmitted to the microprocessor of the electronic control unit of the vehicle, amongst the state and faults signals.\nThis solution fulfils most of the safety requirements inherent in the function, but reveals certain disadvantages.\nIn fact, the information Diag 1. Ond. is made available to the electronic control unit of the vehicle only when the microprocessor to control the electrical machine is functional. This involves a certain software delay in order to become reactive when the information is not present (delay to detect the absence of information).\nAnother disadvantage is that this solution requires use of a microprocessor with a large processing capacity, since the microprocessor to control the electrical machine must control and monitor the machine, as well as support all or part of the system strategy, whilst providing the interface towards the exterior.\nThis solution is also not suitable for the case in which the electrical machine is used both for re-starting of the thermal engine and for torque assistance.\nA simple control line Auth. Ond. cannot be used to prevent both re-starting and torque assistance operations. In fact, the re-starting function can be activated only if the traction chain is open, whereas the torque assistance function can be activated only when the traction chain is closed. These two conditions are contradictory.\nThere is consequently a need for a new architecture which makes it possible to eliminate all of these problems, whilst increasing the level of safety of the global functionality."}
{"text": "This invention relates to golf swing practice devices, and, more particularly, to golf swing practice devices that will assist in preventing undesirable swaying by a golfer during his back swing.\nA primary object of the present invention is to afford a novel golf swing practice device.\nAnother object is to afford a novel golf swing practice device for teaching a golfer to avoid undesirable swaying when he is making his back swing.\nGolf swing practice devices for assisting in training a golfer to avoid undesirable swaying during his back swing have been therefore known in the art. However, golf swing practice devices of the aforementioned type heretofore known in the art have commonly had several inherent disadvantages, such as, for example, being difficult to use; being large and cumbersome in size; being complicated in construction and operation; and being impractical and inefficient in operation, and the like. It is an important object of the present invention to overcome such disadvantages.\nAnother object of the present invention is to afford a novel golf swing practice device of the aforementioned type which may be quickly and easily assembled and disassembled.\nA further object is to afford a novel golf swing practice device of the aforementioned type which, when assembled or disassembled is relatively small and compact in size.\nAnother object is to afford a novel golf swing practice device of the aforementioned type which is simple and uncomplicated in operation, while being highly reliable in operation.\nAnother object of the present invention is to afford a novel golf swing practice device of the aforementioned type which not only affords a golfer practicing therewith an indication of a tendency to sway during his back swing, but, in fact, affords effective protection against such swaying.\nA further object of the present invention is to afford a novel golf swing practice device of the aforementioned type which may be quickly and easily adjusted for use by golfers of different size or build.\nAn object ancillary to the foregoing is to enable such adjustment to be made in a novel and expeditious manner.\nAnother object of the present invention is to afford a novel golf swing practice device of the aforementioned type which is practical and efficient in operation and which may be readily and ecomonically produced commercially.\nOther and further objects of the present invention will be apparent from the following description and claims and are illustrated in the accompanying drawings which, by way of illustration, show a preferred embodiment of the present invention and the principles thereof and what I now consider to be the best mode in which I have contemplated applying these principles. Other embodiments of the invention embodying the same or equivalent principles may be used and structural changes may be made as desired by those skilled in the art without departing from the present invention and the purview of the appended claims."}
{"text": "When dealing with flowing fluids, it is often desirable to reduce the non-uniformity and turbulence in the flowing fluid. Generally speaking, any structure placed in a fluid will cause some turbulence when the fluid begins to flow. In addition, as the speed of a flowing fluid increases, the effects of the turbulence become more pronounced. For example, when testing aircraft in a wind tunnel, it is desirable to produce flowing air having a uniform flow with very low amount of inherent isotropic turbulence. Producing such a low turbulence fluid flow allows the users of the wind tunnel to more precisely monitor the effect of the flowing air in the tunnel on a model placed in the tunnel. If there is a large amount of turbulence in the flowing air in the tunnel, the turbulence will mask some of the effects of the model's structure on the fluid flow and vice versa.\nThe design of a flow straightening and Turbulence Reduction System (TRS) such as a wire-mesh screen for a wind tunnel application is based on a desire to maximize the attenuation of the flow unsteadiness and non-uniformity in the stilling chamber of the wind tunnel. To meet the very stringent flow uniformity and turbulence requirements in the test section where the test subject is positioned, mesh screens are integrated into the stilling chamber design. The stilling chamber is typically designed for operation at relatively lower flow speeds that have sub-critical Reynolds numbers based on the diameter of the fine wire used to create the screens. A Reynolds number is a non-dimensional number used to describe the turbulent nature of a flowing fluid. The Reynolds number of a flowing fluid is calculated as R=ρVD/μ where p is equal to the fluid density, V is equal to the velocity of the flowing fluid, D is equal to the characteristic distance or diameter, and μ is equal to the viscosity of the fluid. At Reynolds numbers larger than about 40 and higher, circular, cylindrical bodies, such as wire, positioned in a cross flow develop unsteady viscous wakes. The viscous wake is convectively unstable and rolls into large-scale coherent rotational flow regions. At higher Reynolds numbers, this is called the Karman vortex street and is composed of positive and negative vorticity that is present in concentrated patterns. These vortices develop higher frequency instabilities and eventually decay into turbulence.\nThe design of a Turbulence Reduction System (TRS) for most wind tunnels is based not only on modifying and reducing turbulence generated in the flow circuit by turbulence producing sources such as fans, heaters, coolers, turning vanes etc., but also on insuring that the turbulence generated by the TRS itself is small or negligible. TRS is generally formed by combination of perforated plates, honeycombs, and screens. Power losses associated with the performance of the TRS are proportional to the square of flow velocity through them. To reduce power consumed by TRS and to meet flow uniformity requirements, the stilling chamber is made as large as practically possible to lower the velocity in the chamber and the screen wire diameter is selected to be very small to achieve sub-critical Reynolds numbers. For large wind tunnels, there are at least two problems associated with this methodology. The first problem is that a wire screen extended over a large span does not retain its intended planar shape. The screen is deformed into a somewhat spherical shape that changes the flow direction in proportion to the local inclination of the screen relative to the flow upstream of the screen. This deformation is the result of a pressure drop that occurs when the fluid flows through the screen. The flow downstream of the deformed screen is, therefore, non-uniform and produces flow angularity in the test section. The second problem is the large required size of the stilling chamber necessary to produce the contraction ratio that is essential to reducing flow stream turbulence close to the desired value.\nOn the basis of both cost and mechanical considerations, the stilling chamber size requirement is especially limiting for the design of high Reynolds' number facilities which are often pressurized. Therefore, in high Reynolds number facilities, the stilling chamber is typically smaller in diameter, has a lower contraction ratio, and has a higher flow velocity for a given test section velocity than a non-pressurized tunnel of the same test section size. This places two contradictory requirements on the screens. The wire diameter must be kept small to operate at sub-critical Reynolds numbers. However, the wire diameter must be increased to reduce deformation due to increased loads on the screen and in turn, increases the turbulence level in the test section and requires a longer, and more expensive, contraction section to provide for increased decay of turbulence exiting the screen.\nThus, the current methods for reducing turbulence in a wind tunnel are to make the walls of the tunnel as smooth as possible and to place a screen across the fluid flow in the wind tunnel. Unfortunately, the smoothness of the walls does not help reduce turbulence in the center of the fluid flow and the screens tend to bend and produce turbulent vortices at higher flow velocities. Therefore, an improved structure and method for reducing turbulence in a wind tunnel are needed.\nThere are many other instances where it is desirable to decrease the turbulence and increase the flow uniformity of a fast flowing fluid. For example, with regard to engines, it is desirable to manage the flow of gas and air into the engine to increase the efficiency and the uniformity of the performance of the engine. This is typically accomplished by injecting the gas into the combustion chamber with a nozzle that creates a fairly uniform mixture of gas and air. In addition, valves are used to precisely control the amounts of gas and air injected into the engine. However, there are currently no methods for reducing the internal turbulence in the flowing fuel and air mixture itself. Turbulence in these mixtures results in small changes in the amount of fuel or air provided to the engine and in the uniformity of fuel air mixture. This variability in flow uniformity reduces the engines' overall performance and efficiency. Therefore, there is a need in the prior art for a method and structure for more effectively reducing turbulence and insuring flow uniformity in an engine."}
{"text": "Dental sealants and adhesives are widely used in clinical settings. Desirable properties include safety, efficacy, durability, and favorable cosmetic properties. It is preferred that dental compositions be shelf stable, easy to formulate, and that they do not set so rapidly as to make them difficult to apply to a patient.\nDental compositions frequently contain monomers which are polymerized by the dentist or technician (e.g. by light, self-cure, or dual-cure). However, many dental compositions form a problematic “oxygen inhibited layer” (OIL) or “uncured layer” on their surface. This layer's polymerization is inhibited due to the presence of molecular oxygen radicals in ambient air. As a result, incomplete polymerization occurs. Such layer often render the surface sticky or tacky, making the dental composition more difficult to mold or shape. Such incomplete polymerization also tends to lead to lower hardness of the surface and/or no curing if a thin surface is present.\nExtoral is a visible-light cured dental resin formulation sold by AFR Imaging Corp. (Portland, Oreg.). Extoral can be used for surface treatment or as a denture resin. Extoral cures rapidly and produces a glossy, hard surface upon irradiation with normal dental light. Extoral is significant in that it does not have an “oxygen inhibited layer”, even when cured with low intensity light. One difficulty with using Extoral is its volatility, leading to a very strong odor. The smell is objectionable to both patients and dentists/technicians, making it difficult to use in a laboratory, and nearly impossible to use in a clinical setting. Another difficulty is the very brittle surface created by Extoral, which may limit its use in certain dental applications such as surfaces subject to compressive forces.\nExtoral's odor is suggestive of the presence of methyl methacrylate. Methyl methacrylate is a volatile compound used in several dental products. In addition to its unpleasant smell, exposure to methyl methacrylate has been linked to various health concerns. Numbness, paraesthesia, reduced pulmonary function, and reduced respiratory function have been observed in dental technicians who have been chronically exposed to methyl methacrylate (Sadoh, D. R. et al., British Dental J. 186(8): 380–381, 1999; Nishiwaki, Y. et al., J. Occup. Health 43: 375–378, 2001). The U.S. Environmental Protection Agency describes methyl methacrylate as an irritant of the nose and throat, and mentions that exposure for short periods of time can cause headache and fatigue (EPA 749-F-95-014 Fact Sheet, November 1994).\nIt would be of great value to develop a dental composition that does not have an “oxygen inhibited layer” and that does not contain methyl methacrylate or other volatile compounds that are irritating and potentially dangerous to dentists, dental technicians, and their patients."}
{"text": "1. Field of the Invention\nThe present invention relates generally to polyaluminum cleaning and laundry compositions. In one embodiment, the present invention relates to polyaluminum cleaning compositions for use on hard surfaces. In one embodiment, the present invention relates to polyaluminum cleaning compositions for use on glass surfaces. In one embodiment, the present invention relates to polyaluminum fabric softeners. In one embodiment, the present invention relates to polyaluminum air/fabric refresher. The invention also relates to cleaning compositions for use with cleaning substrates, cleaning heads, cleaning pads, cleaning sponges and related systems for cleaning hard surfaces.\n2. Description of the Related Art\nPolyaluminum compounds represent a series of water-soluble, inorganic, polymeric compounds with a high charge density and ranging in aluminum oxide content, degree of polymerization and acid neutralization. Due to their high charge density, these materials have typically been used in the water treatment industry for removal of sediment, heavy metals, natural organic matter and microorganisms. Polyaluminum compounds have also been used in the photography and copying industry for improved retention of ink and toner on paper. Polyaluminum compounds have also been used in the pulp and paper industry for improved processing and increased paper strength. Polyaluminum compounds have also been used in the personal care industry for enhanced wetness and odor control for deodorants and antiperspirants.\nHowever, polyaluminum compounds have not been readily utilized in the cleaning and laundry industry. Addition of polyaluminum compounds to cleaning substrates result in improved release of quaternary ammonium biocides as well as improved dirt removal from a surface and retention on a wipe. In fact, polyaluminum compounds can be used in any application where improved quaternary ammonium release or improved soil removal is desired and obtained from the use of cationic organic polymers or chelating agents such as applications related to cleaning hard surfaces. Polyaluminum compounds may be also be used in air and fabric refresher products to reduce the presence of odors and to prevent dust particles and allergens from being released into the air. Additionally, polyaluminum compounds may bind allergens and immobilize microorganisms thereby reducing the spread of illnesses or diseases and consequently, improving the quality of the air being breathed. Since polyaluminum compounds do not contain organic groups (i.e. amines), polyaluminum compounds do not degrade upon heating or upon exposure to UV light and therefore do not discolor fabrics. Therefore, there are some tangible benefits to using polyaluminum compounds in the cleaning and laundry industry.\nPolyaluminum compounds are considered natural ingredients because they are derived from mineral resources as opposed to organic materials derived from petrochemical resources. Because of a desire to use renewable resources, natural based cleaners and laundry products are gaining increasing interest over traditional synthetic products. Most of the cleaners in the market contain only some natural ingredients. One difficulty in formulating natural based cleaners is achieving acceptable consumer performance with a limited number of natural components compared to highly developed formulations using synthetic surfactants and solvents.\nTypical cleaning formulations require multiple surfactants, solvents, and builder combinations to achieve adequate consumer performance. For example, U.S. Pat. No. 5,025,069 to Deguchi et al. discloses alkyl glycoside detergent systems with anionic, amphoteric and nonionic surfactant ingredients. U.S. Pat. No. 7,182,950 to Garti et al. discloses nano-sized concentrates with examples using Tween® surfactants. U.S. Pat. No. 6,831,050 to Murch et al. discloses toxicologically acceptable cleaners containing oleic acid and citric acid. U.S. Pat. No. 6,302,969 to Moster et al. discloses natural cleaners containing anionic surfactants. U.S. Pat. No. 6,420,326 to Maile et al. discloses glass cleaners with ethanol, glycol ethers, and anionic surfactants.\nPrior art compositions do not combine effective cleaning with a minimum number of ingredients, especially with natural ingredients. It is therefore an object of the present invention to provide a cleaning composition that overcomes the disadvantages and shortcomings associated with prior art cleaning compositions."}
{"text": "1. Technical Field\nThe present invention relates generally to a method and apparatus for preventing illegal leakage of a printout of an electronic document and, more particularly, to a method and apparatus for preventing illegal leakage of a printout by outputting the content and security identifier of an electronic document on paper including a security substance, sensing the security substance, and recognizing the content of the security identifier.\n2. Description of the Related Art\nWith improvement in technology of enterprises, the leakage of advanced industrial technology has steadily increased. Such a technology leakage is mostly caused by an insider who intelligently abuses holes in an already deployed security system. Therefore, a technology for preventing external leakage of output confidential documents and drawings is necessary.\nConventional technologies for preventing leakage of a printout have proposed the use of security printing paper. For example, in Korean Patent Application Publication No. 10-2008-0107977 entitled “Security Printing Paper and Method for Manufacturing the Same”, a detection tag is inserted into printing paper such that the detection tag can be sensed when a printout, which is output on the printing paper, passes through a detection gateway. However, it is difficult to detect a person who actually outputs the printout, and it is also difficult to distinguish between a case in which the printout should be exported and a case in which the printout should not be exported.\nAnother technology disclosed in Korean Patent Application Publication No. 10-2011-0126953 entitled “Apparatus and Method for Supplementing Distribution of Internal Document” provides a method of tracing the illegal use of an original internal document by inserting a watermark into the original internal document. However, it is difficult to detect external export of the internal document into which the watermark is inserted.\nSuch a conventional technology for preventing leakage of a printout has limits in approving export, preventing leakage, and recording relevant information depending on whether or not the external export of an actual printout is permitted."}
{"text": "1. Field of the Invention\nThe present invention relates to an ultrasonic probe, and, more particularly, to an ultrasonic probe which has a weighted piezoelectric unit for generating ultrasonic waves.\n2. Description of the Related Art\nAn ultrasonic probe is utilized as an ultrasonic transducer, which transmits and receives ultrasonic waves, in a medical ultrasonic diagnosis apparatus for providing information on the interior of a disease in a living body, by way of example. The ultrasonic probes are classified into several types by their functions and shape, one of which is an array type ultrasonic probe. The array type ultrasonic probe comprises a plurality of strip-shaped piezoelectric units arranged in their width direction, to which an electronic scanning method such as a sector driving method or the like is applied.\nFIG. 1 illustrates the configuration of a conventional array type ultrasonic probe. The probe comprises a plurality of strip-shaped piezoelectric units 2, each of which has electrodes 1a and 1b formed on both main surfaces, respectively. Piezoelectric units 2 are arranged on backing material 3 in the width direction of piezoelectric units 2. The width direction of piezoelectric units 2 is defined in a left-to-right direction in FIG. 1, while the longitudinal direction of piezoelectric units 2 in a direction from the plane of the drawing sheet to the back. Acoustic matching layer 4 is provided to cover transmitting and receiving surfaces of respective piezoelectric units 2. In addition, filler 5 is filled in gaps between respective piezoelectric units 2. Actually, acoustic matching layer 4 is provided on a piezoelectric plate before it is separated into respective piezoelectric units 2, and then filler 5 is filled, so that filler 5 also extends through acoustic matching layer 4.\nIn regard to the array type ultrasonic probe as described above, JP, 5-23331, A and JP, 7-274292, A propose that each piezoelectric unit 2 is weighted in the longitudinal direction of piezoelectric unit 2 to enhance the vibration amplitude in a central region in the longitudinal direction and to reduce the vibration amplitude in both end regions, thereby providing ultrasonic beams with smaller side lobes.\nFIGS. 2 and 3 are side cross-sectional views each illustrating an example of a weighted ultrasonic probe. FIGS. 2 and 3 illustrate cross-sectional structures in a plane perpendicular to the cross-section in FIG. 1, so that the left-to-right direction in the figures corresponds to the longitudinal direction of piezoelectric unit 2.\nIn the ultrasonic probe illustrated in FIG. 2, piezoelectric unit 2 is divided into a plurality of piezoelectric elements in the longitudinal direction. Here, piezoelectric unit 2 is divided into five piezoelectric elements 2d, 2b, 2a, 2c, and 2e, as illustrated in left to right, each of which is weighted. These piezoelectric elements 2a to 2e are electrically connected in parallel with one another. Piezoelectric elements 2a to 2e are weighted such that the degree of polarization is maximized in central piezoelectric element 2a, and becomes gradually smaller toward piezoelectric elements 2d and 2e, respectively, upon the poling process for piezoelectric unit 2. By thus controlling the degree of polarization in piezoelectric elements 2a to 2e, the vibration amplitude is maximized in central piezoelectric element 2a and minimized in piezoelectric elements 2d and 2e at both ends, even if respective piezoelectric elements 2a are applied with the same driving voltage, resulting in a smaller beam width of ultrasonic waves in the longitudinal direction of piezoelectric unit 2 to provide an ultrasonic beam with suppressed side lobes.\nLikewise, in the ultrasonic probe illustrated in FIG. 3, piezoelectric unit 2 is divided into a plurality of piezoelectric elements, here, 2d, 2b, 2a, 2c, and 2e. Central piezoelectric element 2a is directly applied with a driving voltage, while the remaining piezoelectric elements 2b to 2e are applied with the driving voltage through respective resistors 6. Here, the resistances of resistors 6 are controlled to supply a maximum current to central piezoelectric element 2a and a minimum current to piezoelectric elements 2d and 2e at both ends. In this way, the vibration amplitude is maximized in central piezoelectric element 2a and minimized in both end piezoelectric elements 2d and 2e, thereby providing an ultrasonic beam with suppressed side lobes.\nIn the foregoing conventional ultrasonic probes, however, the piezoelectric unit must be mechanically divided into a plurality of piezoelectric elements in the longitudinal direction for weighting the piezoelectric unit, giving rise to a problem of complicated manufacturing process. Since a signal line is routed for each piezoelectric element to apply the driving voltage, the resulting ultrasonic probe also has a complicated structure. The complexity of the structure as well as manufacturing process causes a reduction in productivity of the ultrasonic probe. Also, the weighting as mentioned above can only change the vibration amplitude in steps, i.e., on an element-by-element basis, thus failing to sufficiently suppress side lobes. Furthermore, if the piezoelectric elements are weighted by changing the degree of polarization for each piezoelectric element, the ultrasonic probe will experience difficulties in controlling even the vibration amplitude."}
{"text": "The invention relates to a lower portion of an electric and electronic enclosure, comprising a bottom and fastening arrangements for fastening a cover to the lower portion with a bayonet-like fastening, and the fastening arrangements comprise a bore and a shoulder-like stopper for receiving a mounting bracket of a spindle-like element.\nThe invention also relates to an electric and electronic enclosure having a lower portion with a bottom, a cover, and fastening means for fastening the cover to the lower portion with a bayonet-like fastening, the fastening means comprising a bore and a shoulder-like stopper arranged to the lower portion, a through-hole, and a stop face formed on the cover, and a spindle-like element that comprises a top end having a shoulder for resting against the stop face of the cover, and at a distance from the top end a mounting bracket for fastening at the shoulder-like stopper of the lower portion to provide a bayonet-like fastening.\nElectric and electronic enclosures of the mentioned type and their lower portions are well known. An advantage of a bayonet-like fastening is that the cover of the enclosure can be quickly and easily fastened to the lower portion of the enclosure. A bayonet-like fastening also has the advantage that it can easily be accomplished in such a manner that the cover can flexibly give way when a possible electric arc ignites inside the enclosure, thus preventing any great damage. One drawback with the bayonet-like fastening is that it is not suitable for situations in which the cover of the enclosure needs to remain fixed to the lower portion. In such a case, ordinary screws are usually used to fasten the cover to the lower portion. In addition, the bayonet screw of the bayonet-like fastening has a complex structure, which makes it difficult and expensive to manufacture in comparison with a conventional screw. A further drawback with a bayonet screw is that it cannot be machine-mounted on the enclosure.\nDepending on whether the enclosure needs to be such that the cover is fastened to the lower portion with a bayonet-like fastening or with ordinary screws, the user has to select the lower portion of the enclosure to suit the end use. Thus, there must be two different lower portions depending on the manner in which the cover is to be fastened: a first type of the lower portion comprises bores with a shoulder-like stopper for receiving the mounting bracket of the bayonet screw, whereas a second type of lower portion comprises a bore with threads for receiving the threads of an ordinary screw. The need to manufacture and store lower portions of two different types is naturally a disadvantageous solution that correspondingly increases costs. Sometimes, it is necessary to change the fastening method of the enclosure cover (from a bayonet-like fastening to a screw fastening or vice versa) as the application of the enclosure changes."}
{"text": "The invention relates to a method of manufacturing acicular goethite by forming a suspension of a poorly water-soluble iron (II) compound in a liquid consisting essentially of water oxidizing the iron to form goethite crystals, and filtering off, washing and drying the geothite crystals.\nOf the four known crystalline iron (III) oxide hydrate modifications (.alpha.-, .beta.-, .gamma.-, .delta.-) of the chemical formula FeOOH, goethite (.alpha.-FeOOH) is the most stable and least soluble phase (solubility product [Fe.sup.+++ ].[OH.sup.- ].sup.3 =10.sup.-44 mol.sup.4 1.sup.-4). For structural reasons the crystallite formation of the goethite is acicular in the direction of the c-axis. .alpha.-FeOOH serves as a starting material for the manufacture of magnetic powders. Such magnetic powders are used in the manufacture of carriers for magnetic recording and reproducing of data, audio and video signals. The magnetic powders may consist of acicular maghemite, .gamma.-Fe.sub.2 O.sub.3, or of acicular iron. They are manufactured by reducing acicular geothite pseudomorphously (that is without the particles loosing their acicularity) to magnetite, Fe.sub.3 O.sub.4, and then either oxidizing the magnetite to acicular .gamma.-Fe.sub.2 O.sub.3 or further reducing it into acicular iron.\nU.S. Pat. No. 3,288,563 (corresponding to German Offenlegungsschrift No. 1447134) describes a method of manufacturing a powder for magnetic recording. The powder consists of magnetic iron oxide, .gamma.-Fe.sub.2 O.sub.3. In the method, an iron (II) compound having a low water-solubility is suspended in an alkaline liquid consisting essentially of water with an (OH).sup.- ion concentration of more than 0.2 normal. The iron is oxidized, by introducing air or another oxidizing gas mixture, to form crystalline iron (III) oxide hydrate with a chemical composition of the formula Fe.sub.2 O.sub.3.n H.sub.2 O, in which 1.ltoreq.n.ltoreq.2. The iron (III) oxide hydrate produced is converted into magnetic iron oxide, .gamma.-Fe.sub.2 O.sub.3, by reduction and subsequent oxidation.\nIn the manufacture of acicular .alpha.-FeOOH as a starting material for magnetic recording materials, the starting materials generally are (strongly) alkaline solutions. For example, air or oxygen is blown into a suspension of Fe(OH).sub.2 gel in 1 to 2.5 molar NaOH, for example, at temperatures between 0.degree. and 60.degree. C. to produce the desired acicular .alpha.-FeOOH.\nIn the manufacture of goethite on a production scale, it is difficult to control particle growth and spontaneous nucleation. This is due to the very poor solubility of the .alpha.-FeOOH (one liter comprises in equilibrium at p.sub.H =14 only 5.10.sup.-11 g Fe.sup.3+ besides solid .alpha.-FeOOH). Accordingly, in the manufacture of goethite on a production scale an irreproducible, rather wide spectrum of particles sizes is found.\nIn order to produce magnetic powders having a coercive field strength as large as possible, the goethite particles should have a uniform structure and a pronounced acicularity. Moreover, in order to achieve an acceptable signal-to-noise ratio the particle sizes should fall within a narrow distribution. Particle assemblies manufactured to have optimum characteristics with respect to these criteria, are said to be thermodynamically \"more monoenergetic\" than particle assemblies manufactured without taking these criteria into account."}
{"text": "Services that provide users with guidance information on public transportation services in the form of automated sounds are in wide use. For example, announcement systems used in buses provide passengers with guidance information by way of voice guidance on, for example, bus stops, bus fares, and the like, which guidance voice is played at an appropriate timing, for example at each of bus stops, upon operation of an operator, with such voice guidance being pre-prepared for each of bus routes. Patent Document 1 discloses a configuration in which information such as names of bus stops is announced to passengers through the generation of voice signals that correspond to voice data based on operations such as an input operation into a device installed in a bus carriage. Patent Document 2 discloses a configuration for the creation of voice data corresponding to a guidance voice whereby, supposing that a content of voice guidance is a sentence such as “the next stop is [name of a bus stop]”, names of bus stops are created by utilizing low-compression coding, while common words such as “next”, “is”, and the like are created by utilizing high-compression coding."}
{"text": "1. Field of the Technology\nThe present technology is directed to automated and dynamic repurposing of large amounts of digital graphic data or files.\n2. Description of the Related Art\nGraphical file information has become a major professional digital tool in business. Graphical file information can consist of items such as images, animations, live video, and CAD drawings. Graphical information is generally stored in vector graphics or raster graphics.\nTraditionally, 2D and 3D graphical information has resided in pockets within an organization in different formats and physical locations which are often inaccessible to various different departments within the organization. Products used to create the different graphical files, marketing and other company departments. The costs associated with creating downstream visual communications for documentation, training and support of complex tasks in the product lifecycle are skyrocketing as products become more complex and varied. The need to gain control of these processes in a coordinated manner is now a paramount concern for any product-oriented enterprise wishing to stay competitive.\nThe efficient and effective use and re-use of engineering and design assets can provide significant business and strategic advantage to any engineering or manufacturing enterprise. In particular, the ability to collate data from different graphical applications, used by, for example, different departments, and output new views of the separate or combined data, would be particularly useful."}
{"text": "There are many computer operating systems in use today and each has its advantages and disadvantages. For example, there are operating systems available from Microsoft Corporation, Apple Inc., and Linux as well as many others. Certain operating systems may have performance issues, such as executing more slowly in certain circumstances.\nFor example, it has been noted that the Mac OS X operating system can be slow in launching large applications such as Xcode, Microsoft Office, iMovie, etc. The process of scanning a computer system for malicious software (i.e., malware) using antivirus software can also be very time-consuming, especially for operating systems such as Mac OS X. A manual scan that takes on order of hours will not be acceptable to many computer users.\nWhile there are many factors that will affect how long it takes antivirus software to scan a particular computer system, it is noted that a file system under an Apple operating system is different than a file system under a Microsoft operating system. For example, the number of actual computer files present in an Apple operating system may be much larger than those present in a Microsoft operating system. It has been estimated that a computer executing Mac OS X may have anywhere from 1,000,000 up to 10,000,000 computer files, while a computer executing a Microsoft operating system may have well under 1,000,000 computer files. In particular, while the browser Safari has more than 4,000 files under Mac OS X, it has only approximately 300 files under the Microsoft operating system.\nIt is also noted that the performance of a typical scan engine used to scan a computer system for malicious software is related to the number of files to scan but not necessarily to the file size. Given that performance of a scan engine may be adversely affected by a large number of computer files to scan, it is desirable to improve the performance of scan engine for any computer operating system, especially those having a relatively large number of files to scan."}
{"text": "Brushless DC motors are now commonplace and are typically driven by signals supplied to the stator windings of the motor in accordance with the outputs of one or more Hall-effect sensors which detect the alignment therewith of the rotor poles of the motor.\nThus the requirement for brushes to effect commutation of the driving signals is obviated.\nHowever, at low rotation speeds, smooth and stable control of the motor becomes difficult due the substantial time delays between successive changes of state of the Hall-effect sensor outputs.\nWe have now devised an arrangement which overcomes the limitations of existing control systems for driving DC brushless motors.\nIn accordance with the present invention, there is provided a control system for driving a brushless direct current (DC) motor, the control system comprising means for providing a measurement of the rotational speed of the motor by combining a first signal representing the rotational speed of the motor derived from a measurement of the back-EMF developed in the stator windings of the motor and a second signal representing the rotational speed of the motor derived from the output of means which sense the position of the rotor of the motor, such that as the rotational speed of the motor is increased from a first rotational speed to a second rotational speed, the combined speed measurement is determined to a greater extent by said second speed signal and to a lesser extent by said first speed signal, and means for supplying driving signals to the stator windings in accordance with the combined speed measurement and the position of the rotor of the motor.\nPreferably at said first rotational speed the combined speed measurement is determined from said first speed signal only and preferably at said second rotational speed the combined speed measurement is determined from said second speed signal only.\nWhilst the respective proportional contributions made to the combined speed measurement by each of said first and second speed signals may vary in one or more discrete steps between said first and second speeds, the respective proportional contributions preferably vary continuously with increasing or decreasing motor speed, thereby avoiding any abrupt transitions in the combined speed measurement.\nMost preferably, within a region between said first and second speeds, the respective proportional contributions made to the combined speed measurement by each of said first and second speed signals vary linearly with respect to one another.\nPreferably the position sensing means comprise one or more position sensing transducers. Most preferably the or each position sensing transducer comprises a Hall-effect sensor for detecting the alignment therewith of each of the rotor poles.\nAlso in accordance with the present invention, there is provided a method of driving a brushless direct current (DC) motor, wherein driving signals are supplied to the stator windings of the motor in accordance with a measurement of the rotational speed of the motor and the position of the rotor of the motor, the combined speed measurement being determined by combining a first signal representing the rotational speed of the motor derived from a measure of the back-EMF developed in the stator windings and a second signal representing the rotational speed of the motor derived from the output of means which sense the position of the rotor of the motor, such that as the rotational speed of the motor is increased from a first rotational speed to a second rotational speed, the combined speed measurement is determined to a greater extent by said second speed signal and to a lesser extent by said first speed signal.\nAn embodiment of the present invention will now be described by way of an example only and with reference to the accompanying drawings."}
{"text": "Directional drilling involves controlling the direction of a borehole as it is being drilled. Since boreholes are drilled in three-dimensional space, the direction of a borehole includes both its inclination relative to vertical (dip) as well as its azimuth. Usually the goal of directional drilling is to reach a target subterranean destination with the drill string, typically a potential hydrocarbon producing formation.\nIn order to optimize the drilling operation, it is often desirable to be provided with information concerning the environmental conditions of the surrounding formation being drilled and information concerning the operational and directional parameters of the downhole motor drilling assembly including the drilling bit. For instance, it is often necessary to adjust the direction of the borehole while directional drilling, either to accommodate a planned change in direction or to compensate for unintended and unwanted deflection of the borehole. In addition, it is desirable that information concerning the environmental, directional and operational parameters of the drilling operation be provided to the operator on a real time basis. The ability to obtain real time data measurements while drilling permits a relatively more economical and more efficient drilling operation.\nFor example, the performance of the downhole motor drilling assembly, and in particular the downhole motor, and the life of the downhole motor may be optimized by the real time transmission of the temperature of the downhole motor bearings or the rotations per minute of the drive shaft of the motor. Similarly, the drilling operation itself may be optimized by the real time transmission of environmental or borehole conditions such as the measurement of natural gamma rays, borehole inclination, and borehole pressure, resistivity of the formation and weight on bit. Real time transmission of this information permits real time adjustments in the operating parameters of the downhole motor drilling assembly and real time adjustments to the drilling operation itself.\nAccordingly, various measurement-while-drilling (MWD) systems have been developed that permit downhole sensors to measure real time drilling parameters and to transmit the resulting information or data to the surface substantially instantaneously with the measurements. For instance, MWD mud pulse telemetry systems transmit signals from an associated downhole sensor to the surface through the drilling mud in the drill string. More particularly, pressure or acoustic pulses, modulated with the sensed information from the downhole sensor, are applied to the mud column and are received and demodulated at the surface. The downhole sensor may include various sensors such as gamma ray, resistivity, porosity or temperature sensors for measuring formation characteristics or other downhole parameters. In addition, the downhole sensor may include one or more magnetometers, accelerometers or other sensors for measuring the direction or inclination of the borehole, weight-on-bit or other drilling parameters.\nTypically, MWD systems, such as the MWD mud pulse telemetry system, are located above the downhole motor drilling assembly. For instance, when used with a downhole motor, the MWD mud pulse telemetry system is typically located above the motor so that it is spaced a substantial distance from the drilling bit in order to protect or shield the electronic components of the MWD system from the effects of any vibration or centrifugal forces emanating from the drilling bit. Further, the downhole sensors associated with the MWD system are typically placed in a non-magnetic environment by utilizing Monel collars in the drill string below the MWD system.\nThus, the MWD system may be located a significant distance from the drilling bit. As a result, the environmental information measured by the MWD system may not necessary correlate with the actual conditions surrounding the drilling bit. Rather, the MWD system is responding to conditions that are substantially spaced from the drilling bit. For instance, a conventional MWD system may have a depth lag of up to or greater than 60 feet. As a result of this information delay, it is possible to drill completely through a potential hydrocarbon producing formation before detecting its presence, requiring costly corrective procedures.\nIn response to this undesirable information delay or depth lag, various near bit sensor systems or packages have been developed which are designed to be placed adjacent or near the drilling bit. The near bit system permits the detection of the potential hydrocarbon producing formation almost immediately upon its penetration, minimizing the need for unnecessary drilling and service costs. The drilling operation, including the trajectory of the drilling bit, may then be adjusted in response to the sensed information. However, in order to use a near bit sensor system and permit real time monitoring and adjustment of drilling parameters, a system or method must be provided for transmitting the measured data or sensed information from the downhole sensor either directly to the surface or to a further MWD system for subsequent transmission to the surface. Various attempts have been made in the prior art to transmit the information directly or indirectly to the surface. However, none of these attempts have provided a fully satisfactory solution.\nVarious systems have been developed for communicating or transmitting the information directly to the surface through an electrical line, wireline or cable to the surface. These hard-wire connectors provide a hard-wire connection from the drilling bit to the surface, which has a number of advantages. For instance, these connections typically permit data transmission at a relatively high rate and permit two-way or bidirectional communication. However, these systems also have several disadvantages.\nFirst, a wireline or cable must be installed in or otherwise attached or connected to the drill string. This wireline or cable is subject to wear and tear during use of the system and thus, may be prone to damage or even destruction during normal drilling operations. For instance, the downhole motor drilling assembly may not be particularly suited to accommodate such wirelines running through the motor, with the result that the wireline sensors may need to be spaced a significant distance from drilling bit. Further, the wireline may be exposed to excessive stresses at the point of connection between the sections of drill pipe comprising the drill string. As a result, the system may be somewhat unreliable and prone to failure, which may result in costly inspection, servicing and replacement of the wireline. In addition, the presence of the wireline or cable may require a change in the usual drilling equipment and operational procedures. The downhole motor drilling assembly may need to be particularly designed to accommodate the wireline. As well, the wireline may need to be withdrawn and replaced each time a joint of pipe is added to the drill string. These disadvantages result in a relatively complex and unreliable system for transmitting the sensed information.\nSystems have also been developed for the transmission of acoustic or seismic signals or waves through the drill string or surrounding formation. A downhole acoustic or seismic generator generates the acoustic or seismic signals. However, a relatively large amount of power is typically required downhole in order to generate a sufficient signal such that it is detectable at the surface. To generate a sufficient signal, the necessary power may be supplied to the generator by a hard wire connection from the surface to the downhole generator. Alternately, a relatively large power source must be provided downhole.\nU.S. Pat. No. 5,163,521 issued Nov. 17, 1992 to Pustanyk, et al., U.S. Pat. No. 5,410,303 issued Apr. 25, 1995 to Comeau, et al., and U.S. Pat. No. 5,602,541 issued Feb. 11, 1997 to Comeau, et al. all describe a MWD tool, a downhole motor having a bearing assembly and a drilling bit. A sensor and a transmitter are provided in a sealed cavity within the housing of the downhole motor bearing assembly, adjacent the drilling bit. A signal from the sensor is transmitted by the transmitter to a receiver in the MWD tool. The MWD tool then transmits the information to the surface. The signals are transmitted from the transmitter to the receiver by a wireless system. Specifically, the information is transmitted by frequency modulated acoustic signals indicative of the sensed information. Preferably, the transmitted signals are acoustic signals having a frequency in the range of from 500 to 2,000 Hz. However, alternatively, radio frequency signals of up to 3,000 mega-Hz may be used.\nFurther systems have been developed which require the transmission of electromagnetic signals through the surrounding formation. Electromagnetic transmission of the sensed information often involves the use of a toroid positioned adjacent the drilling bit for generation of an electromagnetic wave through the formation. Specifically, a primary winding, carrying the sensed information, is wrapped around the toroid and the drill string forms a secondary winding. A receiver may be either connected to the ground at the surface for detecting the electromagnetic wave or may be associated with the drill string at a position uphole from the transmitter.\nGenerally speaking, as with acoustic and seismic signal transmission, the transmission of electromagnetic signals through the formation typically requires a relatively large amount of power, particularly where the electromagnetic signal must be detectable at the surface. Further, attenuation of the electromagnetic signals as they are transmitted through the formation is increased with an increase in the distance over which the signals must be transmitted, an increase in the data transmission rate and an increase in the electrical resistivity of the formation. The conductivity and the heterogeneity of the surrounding formation may particularly adversely affect the propagation of the electromagnetic radiation through the formation. As well, noise in the drill string, particularly from the downhole motor drilling assembly, may interfere with the detection of the electromagnetic signals.\nThus, as with acoustic and seismic signal transmission, in order to be able to generate a sufficient electromagnetic signal, the necessary power may need to be supplied to a downhole electromagnetic generator by a hard wire connection from the surface. Alternately, a relatively large power source may be provided downhole.\nFinally, when utilizing a toroid for the transmission of the electromagnetic signal, the outer sheath of the drill string must protect the windings of the toroid while still providing structural integrity to the drill string. This is particularly important given the location of the toroid in the drill string since the toroid is often exposed to large mechanical stresses during the drilling operation. Further, in order to avoid short-circuiting of the system or a short circuit turn of the signals through the drill string and in order to enhance the propagation of the electromagnetic radiation through the surrounding formation, an electrical discontinuity is provided in the drill string. The electrical discontinuity typically comprises an insulative gap or insulated zone provided in the drill string. An insulating material comprising a substantial area of the outer sheath or surface of the drill string may provide the insulative gap. For instance, the insulating material may extend for ten to thirty feet along the drill string. Thus, the need for the insulative gap to be incorporated into the drill string may interfere with the structural integrity of the drill string resulting in a weakening of the drill string at the gap. Further, the insulating material provided for the insulative gap may be readily damaged during typical drilling operations.\nVarious attempts have been made in the prior art to address these difficulties or disadvantages associated with electromagnetic transmission systems. However, none of these attempts have provided a fully satisfactory solution.\nU.S. Pat. No. 4,496,174 issued Jan. 29, 1985 to McDonald, et al. and U.S. Pat. No. 4,725,837 issued Feb. 16, 1988 to Rubin discloses an insulated drill collar gap sub-assembly for a toroidal-coupled telemetry system. The sub-assembly provides a dielectric material in the insulative gap, while attempting to enhance the structural integrity of the sub-assembly at the gap. Although the sub-assembly may enhance the structural integrity of the drill string, the system still requires the propagation of the electromagnetic radiation through the formation to the surface. Specifically, electromagnetic waves are launched from a transmitting toroid in the form of electromagnetic waves traveling through the earth. These waves eventually penetrate the earth's surface and are picked up by an uphole receiving system. The uphole receiving system comprises a plurality of radially extending arms of electrical conductors about the drilling platform, which are laid on the ground surface and extend for three to four hundred feet away from the drill site. These receiving arms intercept the electromagnetic waves and send the corresponding signals to a receiver.\nU.S. Pat. No. 4,691,203 issued Sep. 1, 1987 to Rubin, et al. is directed at a downhole telemetry apparatus for transmitting electromagnetic signals to the surface. The apparatus includes a mode transducer designed to avoid the need for a toroidal transformer. The transducer provides a total electrical discontinuity in the drill string so that a potential difference can be produced across adjacent conducting faces of the drill string. Essentially, the adjacent conducting faces of the drill string are separated from each other by a predetermined insulative gap. Insulation around the gap is selected to induce optimum earth currents when the electrical signal is applied across the faces. Once the signal crosses the insulative gap, it is conducted to the surface through an upper portion of the drill string, where it is transferred from the drill string through a wire to an input transformer for a surface receiver. Once flowing through the transformed primary, the signal is transmitted through a wire installed in the ground near the surface. The electrical signal from the wire propagates through the earth back to the downhole sensor unit and finally completes its path into the mode transducer.\nU.S. Pat. No. 5,160,925 issued Nov. 3, 1992 to Dailey, et al. and PCT International Application PCT/US92/03183 published Oct. 29, 1992 as WO 92/18882 are directed at a short hop communication link for a downhole MWD system. The system comprises a sensor module, a control module, a host module and a mud pulsar. The sensor module includes a transmitter for transmitting an electromagnetic signal, indicative of the information measured by the sensor, to the control module and a receiver for receiving commands from the control module. The control module includes a transceiver for transmitting command signals and receiving signals from the sensor module. Further, the control module transmits electrical signals to the host module through a hard wire connection, which similarly connects to the mud pulsar.\nBoth the sensor and control modules include an antenna arrangement through which the electromagnetic signals are sent and received through a short hop communication link. The sensor and control antennas are transformer coupled, insulated gap antennas. More particularly, communication between the sensor and control modules is effected by electromagnetic propagation through the surrounding conductive earth. The signal is impressed across an insulated axial gap in the outer diameter of the drill string, represented by the antennas, either by transformer coupling or by direct drive across a fully insulated gap in the assembly. The electromagnetic wave from the antenna propagates through the surrounding conductive earth, accompanied by a current in the metal drill string. As the formation conductance increases and resistance decreases, the maximum frequency with acceptable attenuation will decrease. Preferably, a frequency in the range of about 100 to 10,000 Hz is used.\nU.S. Pat. No. 5,359,324 issued Oct. 25, 1994 to Clark, et al. and European Patent Specification EP 0 540 425 B1 published Sep. 25, 1996 are directed at an apparatus for determining earth formation resistivity and sending the information to the surface. The apparatus utilizes a first toroidal coil antenna mounted, in an insulating medium, on a drill collar for transmitting and/or receiving modulated information signals which travel through the surrounding earth formation. A second toroidal coil antenna is also mounted, in an insulating medium, on the drill collar for transmitting and/or receiving the modulated information signals to and from the first antenna.\nMore recent approaches have involved the use of special drill pipe equipped with data links. The disadvantages of this method include high cost associated with the special pipes and unreliability of the couplings in the joints.\nOptic fiber has been used to provide a broadband telemetry system. U.S. Pat. No. 6,041,872 teaches an apparatus having a bared optic fiber cable stored in a spool. The spool can be fit into the drill string and thus the cable will not interfere with adding additional pipes. That attempt has failed because the naked optic fiber cannot withstand the harsh drilling environment. U.S. Pat. No. 6,655,453 records another attempt using armored fiber optic cable for telemetry purposes. Because of the limited space for the cable spool inside the drill string, cable diameter must be small in order to cover the entire borehole length. A thin cable, however, usually means a weak cable that may break in the harsh drilling environment.\nAs revealed above, there remains a need in the industry for reliable real time data telemetry systems and methods for communicating information between multiple positions in a wellbore. The proposed systems and methods of the present invention therefore, address the disadvantages or difficulties associated with conventional telemetry systems and methods."}
{"text": "1. Field of the Invention\nThe present invention relates in general to the field of information handling system wireless communication, and more particularly to an antenna for an information handling system multiplexed interface with enhanced extensions.\n2. Description of the Related Art\nAs the value and use of information continues to increase, individuals and businesses seek additional ways to process and store information. One option available to users is information handling systems. An information handling system generally processes, compiles, stores, and/or communicates information or data for business, personal, or other purposes thereby allowing users to take advantage of the value of the information. Because technology and information handling needs and requirements vary between different users or applications, information handling systems may also vary regarding what information is handled, how the information is handled, how much information is processed, stored, or communicated, and how quickly and efficiently the information may be processed, stored, or communicated. The variations in information handling systems allow for information handling systems to be general or configured for a specific user or specific use such as financial transaction processing, airline reservations, enterprise data storage, or global communications. In addition, information handling systems may include a variety of hardware and software components that may be configured to process, store, and communicate information and may include one or more computer systems, data storage systems, and networking systems.\nInformation handling systems, and in particular mobile information handling systems, often use wireless networking interfaces to communicate through wireless networks. The wireless networks support a variety of different protocols over a variety of different frequencies and communication ranges. One example of such wireless networks is the wireless wide area networks (WWAN) used by mobile (cellular) telephones, such as 3G and 4G mobile broadband. WWAN networks typically use licensed frequency bands assigned to particular network providers. Another example of such wireless networks is the networks that operate in the unlicensed frequency bands. For instance wireless local area networks (WLAN) use the 2.4 GHz and 5 GHz frequency bands to support shorter range communications typically restricted to a home, office, or retail “hot spot” location. WLAN networks also support peer-to-peer device connectivity for adjacent devices and peripherals. Other unlicensed networks typically used for short range device and peripheral connectivity include Bluetooth and 60 GHz WiGig. Some wireless network devices only receive data, such as GPS and television receivers. Others operate over extremely short ranges, such as near field communication (NFC) devices.\nOne difficulty faced by information handling system manufacturers is adapting systems to include new wireless devices as such devices become available. For example, future wireless communications systems may include broadband wireless access known as White Space Spectrum, real-time multimedia services, DVB-H/FLO/DMB-T mobile television, 802.20 Flash OFDM and WiMax. Adding additional integrated wireless devices into existing portable information handling systems presents substantial electrical and ID design challenges. Consumers who purchase portable information handling systems generally want the lightest and most streamlined form-factor platforms; however, placing multiple radio and antenna structures in a device typically involves a considerable footprint in order to obtain acceptable performance with each device. Each radio and antenna structure has mechanical constraints for the placement of components, electrical constraints for powering the components and thermal constraints to ensure that the system operates at acceptable temperatures. Radio performance depends upon antenna placement and tuning as well as the material used in the device platform. Further, passing required emissions testing with minimal signal interference and antenna coupling tends to further restrict design alternatives for adding radio and antenna structures. EMI emissions and regulatory compliance varies worldwide, which introduces other design challenges. In addition, end users probably will not desire all possible radio configurations in every system so that different housings may include different structures that introduce difficulties in the manufacture process."}
{"text": "1. Field of the Invention\nThis invention relates to a method for biochemical assistance in the decompression of divers using a breathing mixture of oxygen and a nitrogen or hydrogen gas diluent and a product for accomplishing the decompression. More particularly, this invention relates to a process for assisting in the removal of hydrogen gas (H.sub.2) or nitrogen gas (N.sub.2) from the systems of divers breathing H.sub.2 or N.sub.2 mixtures under hyperbaric conditions and a product for accomplishing decompression assistance. This product and method supplements and accelerates the removal of hydrogen and nitrogen gas that occurs spontaneously during conventional decompression of divers.\n2. Description of the Prior Art\nAs soon as a diver submerges beneath the water, the diver begins to build a decompression debt. The debt is created by a greater amount of gases dissolving in the blood and cells as the greater depth causes an increase in the partial pressure relative to the conditions prevailing in 1 atm air. When the diver begins to return to the surface these gases leave solution and can create health and life threatening bubbles in the blood stream and tissues. The amount of dissolved gas in the blood and other tissues is a factor of time and pressure. The dissolution of this gas is not instantaneous but happens over time so, until equilibrium is reached, a shallow dive for a long period of time can cause the dissolution of a substantially equal amount of gas as a deep dive for a short period of time. The limits to deep diving by humans are set by factors related to the increasing mass of the water column over the diver as the diver descends.\nIn order to ventilate the lungs during a dive, a gas mixture containing oxygen must be freely available to the diver at a pressure equal to that of the surrounding environment. This breathing gas cannot be pure oxygen (O.sub.2). If the diver is to spend more than half an hour at a depth greater than 10 m, a pressure equal to 2 atmospheres, the diver might suffer from the cumulative toxic effects of O.sub.2. At depths between 10 and 60 m, it is customary to dilute the O.sub.2 with N.sub.2 as in air. *At depths greater than 60 m, helium is commonly used as the diluent because N.sub.2 exerts narcotic effects at these pressures. However, at extreme depths in excess of 300 m, the density of a helium-O.sub.2 breathing mixture is high enough to make it difficult for divers to ventilate their lungs comfortably, particularly during exercise. At depths greater than 300 m an H.sub.2 -O.sub.2 gas mixture offers the advantage of lower gas density thereby providing a gas that is easier to breathe. An additional difficulty encountered at extreme depths is that divers often experience nausea, tremors, or seizures that appear to be induced by high pressure affecting nervous conduction, a phenomenon known as High Pressure Neurologic Syndrome (HPNS). FNT *All references to gas mixtures with nitrogen are intended to include air.\nThe diluent gas in the breathing mixture dissolves in the diver's tissues in proportion to the solubility of the gas in tissues and the partial pressure of the diluent gas. As the diver ascends and the partial pressure of the diluent gas decreases this causes a decreasing volume of gas to remain in solution in the tissues. The diver must ascend sufficiently slowly that the excess diluent gas can be eliminated via diffusion across the epidermal tissue, lungs, or other mucous membrane body surfaces, with the residual gas remaining in solution. If this ascent rate is exceeded, gas bubbles may form in the diver. When the gas bubbles are large and numerous, or in particularly vulnerable tissues such as the spinal cord and joints, they may cause a painful and potentially seriously debilitating condition known as decompression sickness (DCS).\nThe only method of decompression in current use is to carefully control the rate of a diver's ascent. The rate of ascent chosen is based on past history of ascent rates with minimal incidence of DCS. The earliest of these ascent charts were created by the United States Navy and are called the Navy diving tables. Recently, probabilistic models have been generated to assist this process of predicting safe ascent rates for a given dive. In the event that a diver experiences symptoms of DCS, the diver must be recompressed until the symptoms are relieved, and then decompression reinitiated more slowly, again according to past experience. Decompression is thus inherently dangerous because the diver's tissues must remain continuously in a supersaturated state in order to eliminate the burden of excess diluent gas. Decompression sickness cannot be predicted or prevented with absolute certainty because of its probabilistic nature and because each diver reacts differently. The same diver may have a different reaction at different times because of wellness factors. Therefore, the decompression rate necessary to prevent DCS for any individual can only be an approximation based on prior general experience because all the risk factors involved in the off-gassing rate for a given person are not and cannot be known.\nDecompression can also be extremely time-consuming; it takes 12 days to safely decompress a diver from a 300 m dive that lasts as little as a few hours. During the ascent, the diver is at high risk of injury just dangling between the surface and the bottom. Thus, a method for shortening decompression would reduce a time of great personal risk to the diver as well as reducing expenses of the dive operation.\nThe concept of biochemical decompression was first proposed by Dr. Lutz Kiesow, who was a leading scientist in Diving Medicine at the Naval Medical Research Institute. Dr. Kiesow proposed using hydrogenase to cause biochemical decompression. According to this concept, a diver would breathe a gas mixture containing H.sub.2 and O.sub.2. The diver would be supplemented in some fashion with a hydrogenase enzyme, which is found in many bacteria. The hydrogenase enzyme would convert gaseous H.sub.2 to some other molecule or molecules. This process would reduce the diver's burden of excess diluent gas as the diver ascended, thereby shortening the time needed to decompress safely. Dr. Kiesow did not specify any particular means of using the hydogenase or a situs for the interaction. Dr. Milton Axley tried to put this concept in action by trying to put the purified hydrogenase into the red blood cells themselves. Dr. Axley found he could encapsulate the enzyme into red blood cells, but could not devise an animal model in which to test the cells. This work is reported in part in an Abstract for the 5th Meeting of the International Society for the Use of Resealed Erythrocytes as ENCAPSULATION OF HYDOGENASE INTO RED BLOOD CELLS FOR THE PURPOSE OF BIOCHEMICAL DECOMPRESSION, Axley, Kayar, & Harabin 1993, article published in Advances in the Biosciences, Vol 92, pp119-124, 1994 (Elsevier). A study on the kinetics of the concept for use in blood was published as KINETIC MECHANISM STUDIES OF THE SOLUBLE HYDROGENASE FROM ALCALHGENES EUTROPHUS H16, by Keefe, Axley, & Harabin, Archives of Biochemistry and Biophysics, Vol 317, No. 2, pp 449-456, Mar. 10, 1995. The concept of incorporating the hydrogenase into the blood was not further studied because even if the enzyme could be packed into the blood cells and be injected into a diver, the cells could only be circulated for a few weeks before the red blood cells died naturally and were eliminated through the spleen. The foreign protein of the injected enzyme could lead to splenic failure. There remains the concept of biochemical decompression, actually assistance to decompression because normal dissolution of gas continues to occur, without a means of effecting the biochemical decompression."}
{"text": "This invention relates to a process for extruding high viscosity polymeric materials. More specifically, it provides a process for making polymeric extrusions from high viscosity polymeric materials without generating surface roughness on the extrudate, and it permits extruding at higher rates and lower temperatures and reduced energy consumption.\nMany polymeric materials are shaped into useful films, rods, tubes or are coated onto wire by extruding them through an orifice of relatively small dimensions. In such extrusion processes, the difficulties in obtaining a smooth extrudate increase as the viscosity of the polymeric material or output increases. The pressure required to force the material through the exit orifice also increases and an unstable flow phenomenon, often called melt fracture, occurs which results in a rough surface on the extrudate. Such rough surfaces can have adverse effects on extruded products, such as polyolefin rubbers, when, for example, the extrudate with a rough surface is pelletized under water. The interstices on the rough surface retain water and this makes drying the polymer difficult or impractical.\nIt is known that extrusion pressure can be reduced by forming the shaped die from a porous metal material and forcing a gaseous fluid through the pores. However, lubrication with a gas does not provide laminar flow or a smooth surface on the extrudate. It is also known to previously coat shaped extruded materials with various liquids using channels made of porous metal. Such methods have been typically utilized to lubricate the final shaping operation in manufacturing pipe or rods wherein the material is pulled through the die, rather than flowing through by upstream pressure. A method for reducing melt fracture of certain polymeric materials is known in which a small portion of the extruded material is separately heated to a higher temperature than the main mass (to reduce its viscosity) and then is applied as a surface layer in the forming die. This method is difficult to control, limited in viscosity difference, and is not as widely applicable as that of the present invention."}
{"text": "The present invention relates to a composite body for long-term delivery of effective substances, such as multi-layer membranes, multi-layer hollow fibers, or porous powder provided with a skin or a membrane layer, etc., which deliver over a long period effective material in identical quantities, such as pharmaceutics, hormones, chemicals, etc.\nNew and improved effective materials are always developed for the pharmaceuticals industry. In contrast, the methods of application of such effective materials have not been changed. The majority of medicaments are always applied orally or, more seldom, under the skin, intramuscularly or intravenously. The dosing of the effective substance for a successive therapy often has the same great importance as the effective substance itself.\nFor eliminating the disadvantages of the conventional application, it has been proposed to deliver the effective substance from a storage with a certain rate. There are, for example, so-called microcapsules in which the effective material, in general a medicament, is surrounded by a membrane. In this case, it is possible to control the discharge speed of the medicament, for example, by the structure of the membrane, its thickness, selection of the polymer for the membrane, and the like. It has also been proposed to load microporous powder with effective substances, such as fertilizers, insecticides, etc., which deliver in the course of time the effective material to its surroundings.\nIn all the above described methods, it cannot be in many cases avoided that at least in the initial period an increased delivery of the effective substance takes place, or a so-called \"burst\" effect. During the course of time, frequently the quantity of delivered effective substance decreases, despite the fact that the storage has not been exhausted.\nIt is also known to use in such processes composite membranes or porous grains which contain an effective material and are provided with an impregnating substance. The European Pat. No. 5,302 discloses, for example, a process in accordance with which the grains are loaded with an effective material, and a polyurethane skin is formed by chemical reaction to close the pores. Such a method is complicated and difficult. Despite the fact that a whole range of bodies for the long-term delivery of effective materials is known, there is a demand for improved bodies for the long-term delivery."}
{"text": "Stormwater retention basins store water accumulated during a rain event and release the water at a controlled rate to prevent or limit downstream flooding and/or limit downstream waterway erosion.\nSedimentation ponds and sedimentation basins (hereinafter referred to as sedimentation basins) are designed to trap stormwater sediment within the boundaries of a property and prevent the sediment from flowing into downstream waterways. However, sedimentation basins are also often also used to meter the flow of stormwater to downstream waterways, thereby reducing the potential for flooding and downstream erosion.\nSome basins have a discharge opening or orifice with a fixed cross sectional flow area located near the bottom of the basin. As the water level in the basin increases during a storm event, the rate of water discharged from the basin (which is essentially proportional to the square root of the water depth) increases with water level. Because the basin may fill quickly during a storm event, water is discharged from the basin at a maximum rate when the basin is at its most full condition (which is normally soon after the storm event has occurred).\nStormwater events often cause downstream flooding and also scour sediment from the bottom and sides of waterways. The greater the flow of water, the worse the problem becomes downstream. Site development, which tends to include impermeable surfaces such as parking lots, roofs, and the like, normally acts to increase the rate of site discharge and contributes to downstream flooding and erosion problems.\nTypically, all of the sources of stormwater for a given waterway are discharging at their maximum rates during or shortly after a rain event. These sources include underground and aboveground stormwater storage systems, conventional stormwater collection systems, and also overland flows (sheet flows). Given that site development increases the potential for downstream flooding and stormwater volume related problems, it is often advantageous to postpone or delay the onset of the maximum rate of discharge from a basin after the storm event, or to limit the maximum rate of discharge from the basin after a storm event.\nFaircloth, Jr. U.S. Pat. No. 5,820,751 discloses a device that maintains a drain inlet orifice at a constant distance below the water surface of a basin despite raising or lowering of the water level in the basin. Because the inlet orifice is maintained at a constant distance below the surface, the water flows into the orifice at essentially a constant flow rate that is also essentially independent of the basin's water level. Even after the basin fills with water during a storm event, water discharges at that constant flow rate. This can be advantageous because it eliminates the early peak flows from the basin and protects downstream waterways from those peaks. This is especially important in that the peak flows from the various sources which contribute to the flow of a given waterway normally are all releasing at their maximum rate at the same time.\nEven with a constant flow device controlling one or more sources, downstream waterways will still be subjected to maximal flows from sources that are not volume/rate controlled. These sources may include conventional stormwater collection systems and normal flows across natural surfaces leading to the streams and rivers.\nThere is a need, therefore, to delay the peak flows from a basin until preferably well after the rain event. The basin would preferably release water at a relatively lower rate early in a storm event near the start of the basin discharge period, and would discharge a relatively higher rate (including the maximum rate) more towards the end of the basin discharge period."}
{"text": "The invention relates to a flap stops device for a main rotor of a rotorcraft such as a helicopter, and deals more particularly with a device with anti-cone stops or upper flap stops which can be retracted in flight, for a main rotor of the hinged type for any helicopter, and in particular for a helicopter of high tonnage, in which each main rotor is equipped with blades of high mass and large dimensions.\nThe flap stops device according to the invention is more specifically, although not exclusively, intended to equip heavy helicopters, in which the blades of each main rotor have to be able to be folded up, and employed in operating conditions causing them to encounter strong and gusting winds, as is the case with heavy helicopters on board surface ships, on the decks of which the blades of the main rotors have to be manoeuvred in terms of folding up and in terms of deployment, automatically or manually, while the ships are in motion.\nIn order to limit the angular flap or downwards bending of the flapping masses, each mainly composed of a blade, of a part for attaching the blade to the hub, such as a sleeve, and of various other components necessary for operation, such as pins and pitch levers, on hinged main rotors of helicopters under the effect of the self weight of the flapping masses, at low rotational speeds of the rotors and when the rotors are stationary, the rotor heads are generally equipped with flap stops known as \"droop restrainer stops\", and many different embodiments of droop restrainer stops devices have already been proposed.\nIn one embodiment which is well known and particularly popular owing to its advantages as regards the simplicity and low cost of the structure, and as regards the high degree of safety and reliability of operation, the device comprises a central droop restrainer stop common to all the blades and including a rigid so-called \"droop restrainer\" ring mounted so that it can move radially with respect to the rotor mast, and against which, when the rotor is stationary or at low rotational speeds of the rotor, the bearing tracks of rigid lower bearing shoes bear, each of which shoes is secured to a lower part of the root of a blade or of a member for linking the blade to the hub of the rotor.\nWhen the rotor is turning on the ground at a speed above a given threshold, corresponding to a speed which is sufficient to allow the rotor to be controlled effectively using the flight controls or, in flight, for example when it is turning at its nominal speed, the laterally shiftable droop restrainer ring allows the disc of the rotor to be inclined using the cyclic pitch control, because the blade entering the lowest position in its rotation about the axis of the rotor pushes the droop restrainer ring back on the opposite side towards a blade entering the highest position, so that the droop restrainer ring can shift laterally without opposing the movements of the blades.\nWith the rotor halted, the droop restrainer ring also allows the blades to be controlled in terms of pitch in order to control the free excursion of the pitch control, because the control of a change in the pitch of the blades bearing on the droop restrainer ring via their lower bearing shoe causes the droop restrainer ring to be driven through a small rotation by the blades with slight gliding of the bearing surfaces.\nHowever, a droop restrainer stops device with a droop restrainer ring as described for example with reference to FIGS. 3 and 5 of Patent FR 2,427,251exhibits a drawback when a gust of wind, for example, arrives during the critical phases of starting-up or stopping the rotor, when its rotational speed is below the aforementioned threshold, and thus when the centrifugal force is therefore not yet high enough to stabilize the blades in their plane of rotation. Under these conditions, what happens is that the gust of wind may push one of the blades hard downwards, so that it presses violently against the droop restrainer ring which, not being subjected to sufficient reaction forces from the other blades, assumes an extreme radially offset position in which it does not prevent the violently pushed-down blade from reaching a dangerous position of great inclination and possibly from striking or chopping off by its end the tail boom or cabin of the helicopter. Likewise, and still when the rotor is stationary or at low speed, a blade may be lifted by a strong gust of wind and assume a large cone angle, then fall back heavily onto the droop restrainer ring, damaging the latter and/or the lower bearing shoe of the blade, which itself is subjected to an instantaneous bending moment which is very much higher than the loads which it can normally withstand, and this may lead to it becoming unserviceable.\nIn general, starting from the position of equilibrium of the flapping masses, in which there is equilibrium between the static moments of the flapping masses and the reactional moments of contact of the droop restrainer stop tracks with the droop restrainer ring, which is free in the radial plane with respect to the axis of the rotor and ensures equilibrium of the whole by being subjected to forces of contact of all of the blades (which forces are opposed in pairs in a four-bladed rotor), if, for whatever reason, one of the flapping masses pivots upwards, equilibrium is lost. The droop restrainer ring which is unloaded on one side is driven onto this side by the opposite flapping mass which, thus freed, droops downwards, pivoting about its axis of flap. This downwards movement of one blade puts the crew, and the helicopter, and nearby equipment in great danger.\nIn addition to the wind, the reasons which may cause the equilibrium of the flapping masses to be lost, are circumstances which induce forces on the flapping masses, such as movements of the ship and/or the handling of the helicopter and/or the folding-up or deploying of the blades. When a blade is folded from the front towards the back of the helicopter, equilibrium is lost when the centre of gravity of the corresponding flapping mass changes to that side of the flap axis of this blade where the rotor axis is located, which reverses the direction of its static moment, and therefore causes the flapping masses to lose equilibrium. The blade opposite the folded-up blade bends rapidly downwards until it is halted by a flight flap stop, which was not designed for this. During deployment of a blade, the loss of equilibrium of the flapping masses causes a tilting which is the opposite of the aforementioned one and just as dangerous.\nIn order to overcome this drawback, without abandoning the use of a droop restrainer stop with droop restrainer ring, Patents FR 2,434,079 and FR 2,636,914 propose a flap stops device combining droop restrainer stops with droop restrainer ring and lower bearing shoes for the blades, and upper stops limiting the upward angular excursions of the blades. In such a device, when a blade bends downwards, when the rotor is stationary or turning at low speed, the lower bearing shoe of this blade radially pushes against the droop restrainer ring which itself pushes against the lower bearing shoe of at least one of the other blades of the rotor, in the raising direction. However, this raising is limited by the upper flap stops which, when the rotor is stationary or turning at low speed, are still in the working position. The droop restrainer stops or lower stops therefore cooperate with the upper stops in such a way as to reduce the risk of damage to the rotor.\nIn FR 2,434,079, the anti-cone stops comprise, for each blade, on the one hand an upper bearing shoe secured to the part for attaching the blade to the hub and, on the other hand, an upper stop carried by the hub and which can be retracted or moved out of the way by the action of centrifugal force to which it is subjected when the rotor is turning with a rotational speed above the given threshold and sufficient to allow the rotor to be controlled by the flight controls of the helicopter. This upper stop is a cranked lever pivoting about a pin parallel to the axis of the rotor and acted upon by a return spring so that in the \"ground\" position, the spring keeps the lever such that its first arm, constituting the bearing stop proper, is in a position for halting the upper bearing shoe of the corresponding blade and so that its second arm, carrying a flyweight sensitive to the action of centrifugal force in order to cause the lever to pivot against the return spring, bears, over part of its length, against a ring supporting the pivot of the crank lever. The two arms of the latter form an angle such that the centrifugal force applied to the flyweight when the rotor is turning tends to overcome the return force of the spring and to cause the lever to pivot in a direction such that its first arm is moved away from the upper bearing shoe facing it on the corresponding blade.\nIn FR 2,636,914, the anti-cone stops also include, for each blade of the rotor, an upper bearing member secured to the part for attaching the blade to the hub, and an independent upper stop which can be retracted automatically by centrifugal force, comprising a lever mounted so that it can tilt on a pivot secured to the hub, and a first arm of which is shaped into a stop finger, while a second arm carries a flyweight and exhibits a bearing part intended to come against a bearing surface secured to the hub, as well as a return spring returning the lever towards a \"ground\" position in which the bearing part of the second arm of the lever rests against the bearing surface of the hub and the stop finger of the lever is directed towards the corresponding upper bearing member, so as to limit the upwards excursions of the blade.\nIn these two known embodiments, the stops carried by the hub are individual and specific to each blade, so the device includes a large number of components. Furthermore, in FR 2,636,914, each pivoting anti-cone stop is mounted inside a sleeve or clevis for attaching the blade to the hub, and this poses problems of accessibility for maintenance.\nIn these two known embodiments, the devices are suitable for the main rotors of helicopters of low or medium tonnage, and exhibit structures with which there are associated risks of breakage of the anti-cone stops and risks of inadvertent retraction when they are mounted on the rotors of heavy helicopters, even if these embodiments remain compatible with the possibilities regarding bulk, complexity, and permissible mass in such rotors.\nIn particular, inadvertent retraction or breakage of an anti-cone stop, which may also be the origin of a loss of equilibrium of the flapping masses, may result not only from the abnormal forces mentioned hereinabove, via the flapping masses, but also from forces such as those transmitted by manoeuvring the stick during manual folding of a blade, or inertial loading on a retraction flyweight in the event of vertical impact, of the step impulse type, when, as in FR 2,636,914, the anti-cone stop pivots about a pin perpendicular to the axis of the rotor."}
{"text": "It is desirable for telecom carriers to minimize the number of radio frequency cables used to communicatively couple base stations to the antennae located atop a cell tower. The radio frequency cables, which typically extend from the bottom to the top of the cell tower, are expensive, costly to install, and increase wind-loading on a tower. As the number of radio frequency cables in a cell tower increases, the cost of installing and maintaining cell tower increases.\nFIG. 2 is a block diagram of a prior art embodiment of a multiplexer 29. As shown in FIG. 2, the current technology to reduce the number of radio frequency cables to one cable 155 used to communicatively couple base stations 28(1-N) to the antennae (not shown) located atop a cell tower (not shown) comprises band pass filters 111(1-N)."}
{"text": "1. Field of the Invention\nThe present invention relates to an airbag apparatus which is installed on a vehicle and is adapted to be deployed and inflated between the vehicle and an occupant thereof for protection of the occupant when an impact of a predetermined value or greater is applied to the vehicle when the vehicle is involved in a collision.\n2. Related Art\nConventionally, as an airbag for the driver's seat, there is proposed an airbag apparatus described in Patent Publication No. WO 2005/044643A, for example. The airbag apparatus described in WO 2005/044643A includes an inflator functioning as a gas generator and an airbag and is designed to be deployed and inflated momentarily by the pressure of a gas injected from the inflator when an impact of a predetermined value or greater is applied to a front end of the vehicle when the vehicle is involved in a collision. Namely, in this airbag apparatus, the airbag is brought into abutment with an annular rim portion of a steering wheel so as to be deployed and inflated between the steering wheel and an occupant (the driver) seated in the driver's seat when the airbag is deployed and inflated. In addition, when the occupant is captured by the airbag to press the bag toward the front of the vehicle due to the vehicle being involved in a collision, the airbag is not only stopped and received but also is given a reaction force directed to the rear of the vehicle by the rim portion, so that the occupant is protected by the airbag.\nIncidentally, in recent years, steer-by-wire type steering systems (steering wheels) which do not use the column shafts draw attentions. In the steering systems like this, there occurs a case where in place of a normal steering wheel in which a rim portion is formed into an annular shape, an irregular steering wheel is used in which part of a normally annular shaped rim portion is open (hereinafter, referred to as a “irregular steering wheel”). In this irregular steering wheel, when looked about in a radial direction from a pad portion where an airbag apparatus is installed, it appears that part of the rim portion is open. In addition, as the irregular steering wheels, an irregular steering wheel is known in which a rim portion is made up of a pair of left and right grip portions which are each formed into a minor arc and the overall shape of the steering wheel including both the grip portions and a connecting portion which connects both the grip portions together is formed into substantially an H-shape as viewed from the top.\nIn the event of the irregular steering wheel like this, however, when an airbag is deployed and inflated between both the grip portions (namely, at an opening portion of the rim portion), the airbag needs to obtain a reaction force (a reaction force directed to the rear of a vehicle) which is given thereto when restraining (protecting) the occupant by any other means than the rim portion."}
{"text": "In the polyurethane field, increased interest is being shown in compounds which can be added to the polyurethane polymers to act as fire retardant agents. Particular interest is being shown in compounds which have functional groups reactive with the polyol or polyisocyanate used in preparing the polyurethane so that the fire retardant agent can be copolymerized into the polymer chain. One such group of compounds of this type are the polyalkylene glycol polyphosphites and phosphonates. In general, these materials are prepared by transesterifying a secondary phosphite with a polyalkylene glycol in the presence of an alkaline catalyst such as sodium phenolate or sodium methylate. However, many of these materials have relatively high acidity causing them to react with and thereby deactivate certain catalyst systems generally used in the formation of polyurethane polymers such, for example, as tertiary amine compounds. To alleviate this problem, the polyalkylene glycol phosphonates have heretofore been reacted with materials such as alkylene oxides in order to reduce the number of acid groups on the phosphorus. However, addition of the alkylene groups onto the phosphorus has decreased the relative flame retardancy of these compounds. Alternatively, secondary polyalkylene glycol phosphites have been reacted with carbon tetrachloride or chloral in order to add flame retardant chlorine atoms to the molecule. However, the phosphonates formed in this manner are still relatively high in acidity.\nTherefore, it is an object of the present invention to produce a class of compounds which are compatible with polyurethane foams, which can be copolymerized therewith, which have a high degree of flame retardancy and which are relatively low in acidity. Various other objects and advantages of this invention will be apparent from a reading of the disclosure which follows hereinafter."}
{"text": "The present invention relates to a high performance textile fabric and to garments produced from such a fabric. More particularly, the invention relates to the manufacture of a flame retardant textile fabric suited for use in producing close-fitting garments, such as undergarments, that come into direct contact with the skin of the wearer and provide a protective function. The textile fabric is also has applicability for use in various non-apparel applications.\nGarments of this type can be used by the military, police, firefighters, and in sporting applications. The garments must be comfortable, breathable and must have good moisture wicking properties so that perspiration is wicked away from the skin. In addition, the fabric must be capable of being produced either as a white fabric or of being dyed in a variety of bright lightfast colors. Another important criterion is that the garment must be flame retardant.\nNomex® fiber produced by DuPont is widely used in flame retardant fabrics because of its inherent flame retardant properties. However, fabrics made from this fiber are uncomfortable in hot environments and next to the skin. Additionally, the fiber is available only in a limited number of producer-dyed colors and has an inherent yellow color.\nThere exists a need for a high performance flame retardant fabric that is hydrophilic, exhibiting good moisture wicking properties, and which is comfortable in direct contact with the skin."}
{"text": "Transmit diversity seeks to vary the transmission path for different aspects of a radio signal. The diversity may be created by sending two signals at different times, on different frequencies, or from different locations. More complex forms of transmit diversity send variations of a single packet more than once so that a receiver may combine the two signals to reconstruct the original packet. In some transmit diversity systems, the two signals both contain all of the information while in other systems, the two signals each contain a different part of the information. Even if all of the bits are not received from one signal or the other, the original packet might be reconstructed using error correction, depuncturing, maximum likelihood sequence estimation or other techniques.\nIn Long Term Evolution (LTE), transmit diversity for some signals is provided by dividing a packet into parts. The parts are each sent in a different time slot and on a different subcarrier. In addition, different antennas may be used for the two slots. This provides time and frequency diversity and options for spatial diversity.\nLTE uses orthogonal frequency division multiplexing (OFDM) in the downlink and discrete Fourier transform (DFT)-spread OFDM in the uplink. The basic LTE downlink physical resource may thus be seen as a time-frequency grid as illustrated in FIG. 1, where each resource element corresponds to one OFDM subcarrier during one OFDM symbol interval. FIG. 1 is a grid diagram of the LTE downlink physical resource (3GPP TS 36.211, Third Generation Partnership Project Technical Specification No. 36.211).\nIn the time domain, LTE downlink transmissions are organized into radio frames of 10 ms, each radio frame consisting of ten equally-sized subframes of length Tsubframe=1 ms. FIG. 2 is a diagram of the LTE time domain structure where time moves from left to right across the frame.\nThe resource allocation in LTE is typically described in terms of resource blocks, where a resource block corresponds to one slot (0.5 ms) in the time domain and 12 contiguous subcarriers in the frequency domain. Resource blocks are numbered in the frequency domain, starting with 0 from one end of the system bandwidth. Downlink (DL) transmissions are dynamically scheduled, i.e., in each subframe the base station transmits control information to the remote terminals to indicate the resource blocks assigned to transmit data to each terminal, in the current downlink subframe. The remote terminals are typically referred to in LTE as user equipment or UE. This control signaling is typically transmitted in the first 1, 2, 3 or 4 OFDM symbols in each subframe. A downlink system with 3 OFDM symbols as control is illustrated in FIG. 3. FIG. 3 is a grid diagram of a DL subframe showing control in the first blocks followed by data traffic, with reference symbols dispersed through the grid.\nLTE uses hybrid-ARQ (HARQ), where ARQ refers to Automatic Repeat Request or Automatic Repeat Query, where, after receiving downlink data in a subframe, the terminal attempts to decode it and reports to the base station whether the decoding was successful (acknowledgement, ACK) or not (negative acknowledgment, NACK). In case of an unsuccessful decoding attempt, the base station may retransmit the erroneous data.\nUplink control signaling from the terminal to the base station in LTE consists of hybrid-ARQ acknowledgements for received downlink data; terminal reports related to the downlink channel conditions, used as assistance for the downlink scheduling; and scheduling requests, indicating that a mobile terminal needs uplink resources for uplink data transmissions.\nIf the mobile terminal has not been assigned an uplink resource for data transmission, the L1/L2 (Layer 1/Layer 2) control information (channel-status reports, hybrid-ARQ acknowledgments, and scheduling requests) is transmitted in uplink resources (resource blocks) specifically assigned for uplink L1/L2 control on Release 8 of the Physical Uplink Control Channel (Rel-8 PUCCH).\nFIG. 4 is a grid diagram of a PUCCH showing resources assigned for a signal uplink control message. As illustrated in FIG. 4, these resources are located at the edges of the total available cell bandwidth. Each such resource consists of 12 “subcarriers” (one resource block) within each of the two slots of an uplink subframe. In order to provide frequency diversity, these frequency resources frequency hop on the slot boundary, i.e. one “resource” consists of 12 subcarriers at the upper part of the spectrum within the first slot of a subframe and an equally sized resource at the lower part of the spectrum during the second slot of the subframe or vice versa. If more resources are needed for the uplink L1/L2 control signaling, e.g. in case of very large overall transmission bandwidth supporting a large number of users, additional resources blocks may be assigned next to the previously assigned resource blocks.\nThe PUCCH resource blocks are located at the edges of the overall available spectrum to maximize the frequency diversity experienced by the control signaling. Assigning uplink resources for the PUCCH at other positions within the spectrum, i.e. not at the edges, might fragment the uplink spectrum, making it impossible to assign very wide transmission bandwidths to a single mobile terminal and still retain the single-carrier property of the uplink transmission.\nWhen carrier aggregation is used in LTE, one uplink carrier is designed to carry the HARQ-ACK/NACK bits for all DL carrier Physical Downlink Shared Channel (PDSCH) transmissions. To enable the possibility to transmit more than four bits of ACK/NACK, PUCCH format 3 may be used. The basis for Format 3 is DFT-pre-coded OFDM, as diagrammed in FIG. 5 described below.\nIf the number of ACK/NACK bits is up to 11, then the multiple ACK/NACK bits (which may also include one or more scheduling request (SR) bits) are Reed-Müller (RM) encoded to form 48 coded bits. The coded bits are then scrambled with cell-specific sequences. 24 bits are transmitted within the first slot and the other 24 bits are transmitted within the second slot. The 24 bits per slot are converted into 12 QPSK symbols, spread across five DFTS (DFT Spread)-OFDM symbols using an orthogonal cover code (OCC).\nEach of the five sets of QPSK symbols are cyclically shifted with a shift that depends on the cell identity number (Cell ID), the slot number, and the symbol number within the slot. The sets are then DFT pre-coded and transmitted within the one resource block (bandwidth) and five DFTS-OFDM symbols (time) within a slot. The spreading sequence (OCC) is specific to each terminal (UE) and enables multiplexing of up to five users within the same resource blocks.\nFor the reference signals, cyclic shifted constant amplitude zero autocorrelation (CAZAC) sequences are used. This is shown as a processing diagram in FIG. 5 in which input bits are encoded, scrambled and modulated, then applied to weighting multiplexers and DFTs. The DFTs are applied to inverse fast Fourier transform (IFFT) blocks as shown. This is the DFTS-OFDM based PUCCH format 3 for a UE supporting more than 4 HARQ bits in normal cyclic prefix (CP) subframes.\nThe bandwidth of one resource block during one subframe is too large for the control signaling needs of a single terminal. Therefore, to efficiently exploit the resources set aside for control signaling, multiple terminals may share the same resource block, by the use of the OCC on the symbols containing data. In this case the reference signals may also be multiplexed among the terminals. This is done by assigning the different terminals different orthogonal phase rotations of a cell-specific length-12 frequency-domain sequence and/or different orthogonal time-domain covers covering the subframes within a slot or subframe.\nIn LTE Release 10 specifications, orthogonal time-domain covers are not used on the reference signals, but may be used to accommodate additional orthogonal reference signals for PUCCH Format 3. For example, a transmit diversity scheme may be used if one reference signal resource per antenna port is used.\nIf the number of ACK/NACK bits exceeds 11, then the bits are split into two parts and two RM encoders are used, one for each part respectively. This is known as the dual-RM code. Up to 20 ACK/NACK bits (plus one SR bit) may therefore be supported by PUCCH Format 3. Each encoder in the dual-RM code outputs 24 bits. The bits from the two encoders are combined and converted to 12 quaternary phase-shift keying (QPSK) symbols per slot. These 12 QPSK symbols are then spread across the five DFTS-OFDM symbols per slot using one out of five orthogonal cover codes, as in the single-RM code case.\nThe encoding and multiplexing for these cases of transmission without transmit diversity are diagrammed in FIGS. 6A and 6B. FIG. 6A shows encoding and multiplexing up to 11 uplink control information (UCI) bits into slot 0 and slot 1. FIG. 6B shows segmenting 12-21 UCI bits into Segment 1 and Segment 2. The segmentation divides the codewords between the two segments, by, for example, alternating odds to the upper path evens to the lower path. These are combined, encoded and mapped to 24 QPSK symbols and then though a splitter to slot 0 and slot 1.\nThe Dual Codeword Combiner applies an operation in which {tilde over (b)}0, {tilde over (b)}1, {tilde over (b)}2, . . . , {tilde over (b)}23 is the output sequence from the first encoder and {tilde over ({tilde over (b)}0, {tilde over ({tilde over (b)}1, {tilde over ({tilde over (b)}2, . . . , {tilde over ({tilde over (b)}23 the output sequence from the second encoder and NscRB=12, the number of subcarriers per resource block.\nThe output bit sequence b0, b1, b2, . . . , bB-1 where B=4·NscRB is obtained by the alternate concatenation of the bit sequences {tilde over (b)}0, {tilde over (b)}1, {tilde over (b)}2, . . . , {tilde over (b)}23 and {tilde over ({tilde over (b)}0, {tilde over ({tilde over (b)}1, {tilde over ({tilde over (b)}2, . . . , {tilde over ({tilde over (b)}23 as follows\n            Set      ⁢                          ⁢      i        ,          j      =      0                  while      ⁢                          ⁢      i        <          4      ·              N        sc        RB                                b        i            =                        b          ~                j              ,                  b                  i          +          1                    =                        b          ~                          j          +          1                                        b                  i          +          2                    =                        b                      ~            ~                          j              ,                  b                  i          +          3                    =                        b                      ~            ~                                    j          +          1                          i    =          i      +      4            j    =          j      +      2            end    ⁢                  ⁢    while  \nThe “map to . . . ” operations apply a cell, slot and symbol specific cyclic shift of the symbols in the time domain to provide inter-cell interference randomization. The DFT processing referred to above is also part of the “map to . . . ” operation.\nIn LTE Release 10, space orthogonal resource transmit diversity (SORTD) is used for PUCCH Format 3 to achieve transmit diversity. Transmission with two antenna ports is supported and the Format 3 encoding and mapping shown in FIG. 5 is repeated for each of the two antennas apart from the cyclic shift of the reference signals and the orthogonal cover codes on the data, which are different to provide orthogonality between signals transmitted on the different antenna ports. The encoding and multiplexing shown in FIGS. 6A and 6B may be adapted to map to two antenna ports and two slots for four paths or to support even more ports and slots."}
{"text": "1. Field of the Invention\nThis invention relates generally to computer graphics and, more specifically, to using computer systems to model hair.\n2. Background Art\nComputers are widely used to model figures, from insects and monsters to human beings. While it is relatively easy to use a computer to create and position the larger, simpler body parts of a figure (such as the head, torso, and limbs), it is much more difficult to create and position the smaller, more detailed body parts such as hair and fur. Hair and fur are especially difficult to model because each strand must be created and positioned separately. Since one figure is often covered by thousands of strands of hair or fur, it is obvious that these strands cannot all be created and positioned manually.\nHuman hair is even more difficult to model than fur because of its length and the variety of possible styles. A real human head contains about 100,000-200,000 strands of hair. And for each head of hair, there are dozens, if not hundreds, of possible hairstyles. Art directors often have particular hairstyles in mind, and computers must be able to generate these hairstyles in order to be useful.\nIn the past, there have been three major approaches to modeling hair: the texture map approach, the cluster approach, and the interpolation approach. In the texture map approach, a user paints areas on a figure to specify different characteristics of the hairs growing out of that area. For example, one map might specify hair length, where the area near the forehead is painted so that the hairs will be short, while the area near the back of the head is painted so that the hairs will be long. Another example is specifying the direction of hair growth, such that hairs in one area grow out perpendicularly to the head, while hairs in another area grow out close to the head. Texture maps have also been used to specify density and shape.\nWhile texture maps may be used to specify a variety of characteristics for a hairstyle, it is difficult to use texture maps in order to create complex hairstyles. Since each texture map addresses only one characteristic at a time for a given hair, multiple texture maps are necessary in order to specify multiple characteristics for a given hair. Specifying and modifying multiple texture maps can become labor intensive, so frequently texture maps are used for only gross control of shape, and the number of texture maps used is limited as much as possible. In addition, texture maps are generally a non-intuitive and inefficient way to style hair, depending on which characteristics the texture maps are specifying. As a result, it is frequently easier to create a complex hairstyle by directly modeling the shape curves of the hairs. These limitations make texture maps most useful for simple hairstyles.\nIn the cluster approach (also known as the wisp approach), the user manually creates and positions hair tubes. These tubes have diameters much larger than the diameters of individual hairs. Thus, fewer tubes are needed to cover an area compared to the number of hairs needed to cover that same area. During rendering, each tube is filled with many hairs, so that the hairs follow the direction of the tube. While the cluster approach does allow for individual hair variations within a given tube, these variations usually rely on random values to modify hairs and thus are not very controllable by users. Thus, the cluster approach is best for straight hair. The other problem with the cluster approach is that the resulting hairstyle always looks like a lot of hairs in tubes, rather than a real hairstyle.\nIn the interpolation approach, the user manually creates and positions several guide hairs. While the appropriate number of guide hairs depends on how much manual control the user needs at this phase, for many hairstyles, the number tends to be about 40-100. These guide hairs, also known as control hairs, possess certain characteristics, such as size and shape, and help define key aspects of the hairstyle. The computer uses interpolation to automatically generate hairs to fill in the gaps between the guide hairs. These interpolated hairs have characteristic values similar to those of the surrounding guide hairs. The actual characteristic values assigned to an interpolated hair are based on the origin of the interpolated hair. Many interpolation methods exist and may be used. In one embodiment, the guide hairs that are closer to the interpolated hair (based on the origins of the guide hairs and the origin of the interpolated hair) affect the characteristic values of the interpolated hair more than the guide hairs that are farther from the interpolated hair. This is known as a radial basis function. For example, if three interpolated hairs are evenly spaced between two guide hairs, the interpolated hairs closest to a guide hair will mostly resemble (via characteristic values) that guide hair, and the interpolated hair in the middle will be a mixture of the two guide hairs.\nThe interpolation approach also enables a user to specify a blending weight for each guide hair. A large blending weight increases the effect of a guide hair on a characteristic value of an interpolated hair, while a small blending weight decreases the effect of a guide hair on a characteristic value of an interpolated hair. The problem with the interpolation approach is that it results in artifacts, such as hairs that stand up straight between two guide hairs that go left and right. Also, the interpolation approach cannot be used for clumpy hairstyles, since interpolation tends to create a very smooth distribution of hairs.\nWhat is needed is a way to model hair that enables users to design complex and realistic hairstyles without requiring each hair to be created and positioned manually."}
{"text": "In related art, an imaging device is assembled as a module in which two chips, which are a CMOS image sensor (CIS) chip and an image processor chip, are respectively mounted on packages.\nAlternatively, the chips may be mounted chip-on-board (COB).\nRecently, there has been a need for reduction in mounting areas and in size for mounting an imaging device on a cellular phone or the like, and a System on Chip (SOC) that integrates the two chips on one chip has been developed (see FIG. 2(A)).\nA process combining a CIS process and a high-speed logic process for integration on one chip, however, results in an increase in the number of processes and high cost, and moreover, it is difficult for such process to produce both analog characteristics and logic characteristics, which may cause degradation in the characteristics of the imaging device.\nIn this regard, methods for assembling the two chips at the chip level while both reducing the size and improving the characteristics are proposed (see Patent Literatures 1 and 2)."}
{"text": "1. Field of the Invention\nThe invention relates to a rotary tubular kiln with a satellite cooler. In such a kiln and cooler combination, the cooling tubes are commonly held by support means that allow for thermal expansion of the cooling tubes. Such support means are located in the discharge region of the rotary tubular kiln at two fixing locations on the kiln casing spaced apart in the axial direction. Each of the cooling tubes communicates at its inlet end by way of an inlet pipe with the interior of the kiln.\nRotary tubular kilns of this type are used in the production of cement clinker. When such a rotary tubular kiln is in service, the kiln casing, cooling tubes and support means are subjected to highly variable stresses due, on the one hand, to their heating, which frequently exhibits extreme local variations, and, on the other hand, to the mechanical load cycling of the cooling tubes by the material being cooled owing to the revolution of the kiln.\n2. Statement of Prior Art\nAttempts were first made to prevent the expensive damage such as cracks and fractures of the support means and of the kiln casing which are caused by the co-operation of these effects, by use of intrinsically rigid retaining or bracing elements for the cooling tubes mounted so as to be tiltable in the axial direction on the kiln, casing and on the cooling tube, which permitted an unobstructed thermal expansion of the cooling tube only in the direction of the kiln's longitudinal axis.\nIn view of the unsatisfactory results, attempts were made to give the cooling tubes also a restricted radial thermal expansion by redesigning the support means.\nIn one known arrangement of this type annular support means are illustrated in West German OS 25 19 458. Such a support means comprises a sheet metal plate attached to the kiln casing, the sheet metal plate having an inner and an outer annular zone in the radial direction, using an interposed elastic bush, if desired. The two annular zones can be joined by slotted radial struts constructed integrally with them, or by plate rings.\nThe cooling tubes are accommodated in the plate rings or in the circular orifices formed by the arcuate radial struts so as to leave an annular gap, whilst they are screwed to the retaining elements, each by two lugs provided mutually opposite in the peripheral direction of the sheet metal plate on the radial struts, the plate rings, or the cooling tube shell itself. Where plate rings are used, they should be fixed, likewise by means of screws, to the inner and the outer annular zone of the sheet metal plate.\nIt is further possible for the inner and outer annular zones of the sheet metal plate to be composed of segment-like elements arranged in staggered relationship in the axial direction, whilst the segments overlap partly in the axial direction.\nLastly, a lamellar arrangement of the sheet metal plates is also possible.\nThe prior art support means described permit an unobstructed axial thermal expansion of the cooling tubes during service. However, their thermal expansion in the radial direction is restricted at the lugs, in which local stress concentrations occur, as they also do in the adjacent region of the cooling tube due to the turning of the latter. Having regard to the cyclic loads which act upon the cooling tubes during the revolution of the kiln and which are produced in constantly changing planes due to sagging of the cooling tubes, an increased danger of fracture occurs at these points of concentrated thermic stresses. The screw joints of the plate rings to the inner and outer annular zone of the sheet metal plate are threatened for similar reasons."}
{"text": "The Internet has been plagued by a variety of security threats over the past several years. The security menace on the Internet is exacerbated as more professionals are getting into this lucrative business. Attacks are also getting more sophisticated, as the attackers are not merely interested in achieving publicity. The low rate TCP Denial of Service (DoS) attack is one such intelligent attack, which was first explained in A. Kuzmanovic and E. Knightly, “Low-Rate TCP—Targeted Denial of Service Attacks (The Shrew vs. the Mice and Elephants),” Proceedings of ACM SIGCOMM 2003, Kalrushe, Germany, August 2003, pp. 75-86. A low rate DoS attack is typically illustrated by a periodic waveform shown in FIG. 1. T is the time period, t is burst period, and R is the burst rate.\nA low rate TCP DoS attack exploits widely implemented minimum RTO (retransmission timeout) property of the TCP protocol. The following characterize a low rate TCP DoS attack:                It sends periodic bursts of packets at one-second intervals.        The burst rate is equal to or greater than the bottleneck capacity.        The burst period is tuned to be equal to the round-trip times of the TCP connections. The burst period determines whether the attack will cause DoS to the TCP connections with small or long round trip times.        The exponential back off algorithm of the TCP's retransmission mechanism is eventually exploited.        \nIn a Reduction of Quality (RoQ) attack, the attacker sends high rate short bursts of the attack traffic at random time periods, thereby forcing the adaptive TCP traffic to back-off due to the temporary congestion caused by the attack bursts. In particular, the periodicity is not well defined in a RoQ attack, thus allowing the attacker to keep the average rate of the attack traffic low to evade the regulation of adaptive queue management (AQM) like random early detection (RED) and RED-PD (Preferential Dropping). RED detects congestion at early stages by monitoring the average queue length. RED-PD is another AQM scheme that regulates the long-lived flows, which occupy most of the bandwidth in the Internet. By sending the attack traffic, the RoQ attack introduces transients and restricts the router queue from reaching the steady state. The open knowledge of these stealthy attacks presses for early fixes. For simplicity, the term “low rate DoS attack” refers to both the low rate TCP DoS and RoQ attacks, unless otherwise stated.\nThe threat imposed by the low rate DoS attacks which use IP address spoofing, and the low rate DoS attack model using botnets is significant. Owing to the open nature of the Internet, IP address spoofing can still evade ingress and egress filtering techniques at many sites. A low rate DoS attack can use IP address spoofing in a variety of ways like random IP address spoofing and continuous IP address spoofing. The use of IP address spoofing most importantly divides the high rate of a single flow during the burst period of the attack among multiple flows with spoofed identities. This way, an attacker can evade detection systems that concentrate on finding anomalous traffic rate. The detection systems that rely on identifying periodicity of the low rate DoS attack in the frequency domain can detect the periodicity, but they fail to filter the attack traffic as it is difficult to know the IP addresses that an attacker will use in the future. This problem is further exacerbated by the use of botnets; a botnet is a network of compromised real hosts across the Internet controlled by a master. As an attacker using botnets can control thousands of hosts, it can easily use these hosts to launch a low rate DoS attack like a low rate DoS attack that uses random or continuous IP address spoofing. Now, with the use of botnets, the IP addresses (of compromised hosts) are not spoofed and so these packets cannot be filtered by spoofing prevention techniques. In fact, these attack packets are similar to the HTTP flows. Thus, the objective is to detect a low rate DoS attack that can deploy different IP address spoofing techniques, and then filter the attack traffic. In a low rate DoS attack, an attacker usually targets the network element which is a router. The low rate DoS attacks launched by using botnets can behave like normal connections and can go undetected by prior systems which can, however, be used to regulate and limit both the legitimate and attack traffic. These attacks are similar to the low rate DoS attacks that use continuous IP address spoofing. The possibility of an attacker using UDP packets to launch a low rate DoS attack always exists, and prior systems do not provide a provision to mitigate such threat from UDP packets. An attacker can use aggregate UDP traffic coming from many networks to launch a low rate DoS attack. There is a need to detect and mitigate these problems and related issues."}
{"text": "Portable draft boxes and bars for dispensing beverages are well known in the art, as exemplified by U.S. Pat. Nos. 6,334,329 issued Jan. 1, 2002 to Weller; 6,595,475 issued Jul. 22, 2003 to Svabek; 5,339,986 issued Aug. 23, 1994 to Mihalich; 1,772,111 issued Aug. 5, 1930 to Rice; 6,481,238 issued Nov. 19, 2002 to Jennings; and 5,915,602 issued Jun. 29, 1999 to Nelson. Many are large and cumbersome. Bag-in-box beverage packages are often produced for providing a large volume, 3 liters or more, of wine in a package that dispenses wine without admitting air. They are more economical than bottles, easier to handle, and environmentally friendly. They have a dispensing outlet at a low position, so that the outlet must be positioned over the edge of a support surface. A drinking receptacle must then be hand held below the outlet for filling. This is awkward when dispensing many portions at a public gathering."}
{"text": "1. Field of the Invention\nThe present invention relates to a recording medium discharge apparatus, and image forming apparatus equipped with a recording medium discharge apparatus. In particular, the invention relates to a recording medium discharge apparatus mounted on an image forming apparatus, such as a printer, to receive, sort and store recording media such as sheets of paper. The invention also relates to an image forming apparatus fitted with such a discharge apparatus, and to an image forming apparatus which can be fitted with such a discharge apparatus.\n2. Description of Related Art\nRecording medium discharge apparatuses has have been known which receive, sort and store recording media such as paper sheets discharged from an image forming apparatus such as a copying machine, a printer, or a facsimile machine. In general, a recording medium discharge apparatus of this type includes bins in which paper sheets can be stacked and a means for conveying to the bins the paper sheets discharged from the image forming apparatus. The conveyed sheets are sorted and stored in the bins.\nA level detector or sensor is provided for each bin to detect it being filled with paper sheets. The detector includes a lever or arm, which may be supported near the paper discharge port of the discharge apparatus. The lever can swing or pivot across the direction in which paper sheets can be discharged. The lever is urged for contact with the top one of the sheets stacked in the bin. The detector also includes a sensor for detecting the position of the lever. The detected lever position determines whether or not the stacked sheets have reached a predetermined height or level.\nEach bin may be used like a mailbox allotted to a specified user. When paper sheets are stacked to a certain height in one of the bins, a user may unfortunately take out the sheets which are to allotted to another user. Finding that the sheets are for another user, this user may put them back in the bin. Once the sheets are taken out, the detector lever of the bin turns downward. If the sheets are returned into the bin, their rear ends push the lever toward the discharge port, preventing the lever from leaving the downward position. As a result, the lever blocks the port, which will then be jammed with the next sheet being discharged."}
{"text": "With rapid growth of the Internet of Things (TOT) technology, more and more sensors are deployed to sense various environmental parameters or transmitted signals. Many of such sensors, especially for remote wireless sensing applications, have a very limited operational lifetime. To collect all of the sensors after a project is complete could be costly, time-consuming, or impossible. However, leaving the sensors in place when they are no longer needed, may pose both environmental and information security risks.\nTherefore, there is a need in the art for a cost-effective approach to manage such sensors after their useful life is over."}
{"text": "Puzzles and games have long been a huge source of amusement, entertainment, and even learning for young and old alike. Items such as jigsaw puzzles, crossword puzzles, Rubik's Cube®, Tetris®, and the like can improve problem solving, increase vocabulary, enhance reasoning skills, and can also be a fun way to socialize and improve fine motor skills, among other recognizable benefits. Many types of puzzles and games can be done alone or with other people, such as to foster cooperative play or even to encourage family time.\nThere are various types of puzzle and game categories, such as card games, board games, word games or puzzles, multi-player, solitaire or single player, video or computer games, and the like. While some puzzles and games are geared towards a younger audience, others are directed at an older audience. Although examples such as the aforementioned Rubik's Cube and Tetris do exist, it is often difficult to create a puzzle or game that combines simple rules with diverse skill levels, while still remaining popular and fun with young, old, simpler, and more sophisticated or smart audiences.\nDespite the wide variety of offerings that are available, many puzzle and game enthusiasts still desire further variety when it comes to their puzzles and games. As such, while existing puzzles and games have worked well over many years, there is always a desire for improvement and new variations. In particular, what is desired are new varieties of puzzles and games that combine simple rules with diverse skill levels, that appeal to a wide range of users, and that still remain popular and fun to play."}
{"text": "1. Field of the Invention\nThe present invention relates to carrier recovery systems and carrier recovery methods.\n2. Description of the Related Art\nIn communication systems using Quadrature Amplitude Modulation (QAM), such as OpenCable and DVB-C communication systems, receivers have to compensate for frequency offset and phase offset by blind methods since a transmitted signal has no pilot and training sequence. The blind methods are named carrier recovery.\nThere are many solutions for carrier recovery. One conventional carrier recovery technique compensates for frequency and phase offsets at the output terminal of the equalizer. This conventional technique provides a wide frequency offset estimation range (such as up to 300 KHz) but has a bad performance at the equalizer since the signal inputted into the equalizer includes the uncompensated frequency offset. Another conventional carrier recovery technique compensates for frequency and phase offsets at the input terminal of the equalizer. This conventional technique provides a good performance equalizer but limits the frequency offset estimation range (only up to 40 KHz for example).\nThus, novel carrier recovery techniques providing good performance equalizer and wide frequency offset estimation range are called for."}
{"text": "The present invention relates to the field of interconnections between dampers and ducts, and more particularly, to interconnections which are to be made between ducts, and fire, smoke, or air control dampers which may be part of an overall fire, smoke and/or air handling system.\nIn order to pass \"UL\" testing procedures, and in order to provide a maximum degree of safety in a given installation, the interconnections between ducts and a damper are critical parts of a fire and/or smoke protection system. If, for example, proper interconnection is not made between a fire damper mounted in a fire wall and its associated duct extending into the adjacent room, premature separation of that duct from the damper may occur, thereby introducing smoke or fire into a room which may otherwise be free from same. On the other hand, if the interconnection between the duct and the damper is too secure, falling debris which strikes the duct on one side of the fire wall may act through the duct on the damper to pull the damper from the opening in which it is disposed, thereby leading to the spread of fire through the fire wall and the loss of fire protection otherwise imparted by the damper.\nFire, smoke, and air control dampers and their associated duct work are, of course, installed at the job site. It is accordingly important, in order to achieve the desired separation characteristics between duct work, which may be supplied from one source, and dampers which may be supplied from a second source to provide adaptors which are nearly universal and which are not prone to improper installation. Since dampers and duct-work may be installed by a contractor who is not particularly skilled or knowledgeable concerning the desirable separation characteristics of a suitable damper-duct work interface, it is important to provide a simple positive reception of the duct by the adaptor which will counteract any tendency of the installation workman towards over-fastening the duct to the damper.\nHeretofore, most damper-duct interfaces have been created using slip joints or brackets which may be riveted to portions of the ducts and dampers, with some attempt being made to establish desirable shearing or release characteristics between the ducts and dampers, as for example, by limiting the number of rivets which may be used between the duct and slip joints. See for example my prior issued U.S. Pat. No. 3,727,663 entitled, \"Device For Mounting Fire Damper\"."}
{"text": "Computer-based presentations are facilitated by a projector connected to a computer to project a computer screen image onto a projection surface, such as a projection screen surface or a wall surface. The projected computer screen image can be the viewable output image of the computer, which is commonly displayed on a computer monitor. Thus, a computer-based presentation may involve displaying what would typically be displayed on a computer monitor. However, the projected computer screen image may be an electronic image that is different than the image displayed on the computer monitor.\nIn order to point to a specific item in a projected computer screen image during a computer-based presentation, a presenter may use any laser pointer to point to that specific item in the projected computer screen image. The laser pointer is used to create a laser spot on the projection surface, which allows viewers to know what item to which the presenter is pointing. However, a laser pointer typically cannot be used to control the computer-based presentation. Rather, the computer keyboard and/or the computer mouse are usually used to control the computer-based presentation, which anchors the presenter to the computer and limits the mobility of the presenter.\nRecently, laser pointer tracking systems have been developed to optically track the laser spot on the projection surface to control the cursor of the computer that is being used to provide the projected computer screen image. The laser pointer used in some of these systems may include buttons or other controls, which emulate computer mouse buttons. Thus, using a laser pointer tracking system, the presenter can freely move about without having to return to the computer to control the computer-based presentation.\nA concern with some conventional laser pointer tracking systems is that the systems must be manually calibrated before use. In particular, the position of the projected computer screen image on the projection surface must be manually entered into the system so that the system can determine the location of a laser spot relative to the original computer screen image to move the cursor to that location. Typically, this involves pointing the laser pointer to each corner of the projected computer screen image on the projection surface during a calibration process, which can be annoying to some users.\nThus, there is a need for a laser pointer tracking system, which can automatically perform a calibration process."}
{"text": "Semiconductor packages widely used in electronics (e.g., personal computers), communication devices, etc. have been more highly functionalized or miniaturized in recent years. Along with this, the high integration or high-density packaging of each component for semiconductor packages has been increasingly accelerated. In the high integration or the high-density packaging, the difference in the coefficient of thermal expansion between a semiconductor device and a printed circuit board for semiconductor plastic packages causes the undesired warpage of semiconductor plastic packages. Various measures against this problem have been attempted.\nOne example of the measures includes reduction in thermal expansion of insulating layers for use in printed circuit boards. This approach is to bring the coefficient of thermal expansion of a printed circuit board close to the coefficient of thermal expansion of a semiconductor device to thereby suppress the warpage (see, for example, Patent Literatures 1 and 2).\nIn addition to the reduction in thermal expansion of printed circuit boards, increase in the rigidity of laminates (high rigidity) or increase in the glass transition temperatures of laminates (high Tg) has been studied as an approach for suppressing the warpage of semiconductor plastic packages (see, for example, Patent Literature 3)."}
{"text": "Electromagnetic shields may be used in some computing devices to protect sensitive circuitry from electromagnetic interference that may disrupt operation of the circuitry or to constrain the electromagnetic interference generated by an especially “noisy” component."}
{"text": "1. Field of the Invention\nThe present invention relates to a magentic recording medium and a process for the preparation of the same. More particularly, the invention relates to a magnetic recording medium having a magnetic recording layer consisting of at least two layers which is improved in frequency characteristics and a process for the preparation of the same.\n2. Description of Prior Art\nMagnetic recording media such as an audio tape and a video tape are desired to be excellent in various properties in recent years. For example, an audio cassette tape is required to show low noise, well-balanced frequency characteristics and high output in all frequency bands. Such requirement for the audio cassette tape tends to increase more and more because various developments such as a development of digital system have been recently made in pursuit of ultra hi-fi audio sources or extremely low noise audio sources as shown in a compact disc. A video tape is also required to show higher video output and lower noise.\nIn order to accomplish the above-mentioned excellent properties of the magnetic recording medium, various methods have been proposed. For example, a method of using iron oxide magnetic particles having a short length in the mean longitudinal length as a magnetic powder of the recording medium is known is lower a noise level. This method is able to lower a noise level of the resulting medium, but unable to give satisfactory frequency characteristics to the medium.\nFor improving the frequency characteristics, a magnetic tape having a magnetic recording layer which consists of two or more layers has been proposed. In the conventional magnetic tape having two or more recording layers, an upper recording layer is made to have a higher coercive force than that of a lower recording layer so as to give excellent characteristics in the high-frequency band to the upper recording layer, and thereby to obtain a high output level in all frequency bands in the resulting medium. However, such conventional tape does not have a uniform reproduction output in all frequency bands because of the distinct separation between the upper and lower recording layers, and for example, the reproduction output of the tape lowers in the mid frequency band or partially decreases in the extreme in the reproduction procedure. Further, the conventional magnetic tape having plural recording layers is disadvantageous from the viewpoint of practical operation and cost for the preparation (for example, a number of coating procedures are required for the tape), although the frequency characteristics are enhanced in accordance with the number of recording layers.\nIn other conventional magnetic recording medium having plural recording layers, magnetic particles having a short length in the longitudinal direction are incorporated into the upper recording layer for the purpose of reducing a noise level and magnetic particles having a long length in the longitudinal direction are incorporated into in the lower recording layer for the purpose of increasing frequency characteristics in the high-frequency band, in addition to the above-mentioned consideration on the coercive force. However, when the magnetic particles having a long length in the longitudinal direction is used for the lower recording layer, the surface smoothness of the resulting lower recording layer is decreased more than the case of using magnetic particles having a short length in the longitudinal direction. Thus decreased surface condition of the lower recording layer influences on the upper recording layer, and as a result, satisfactory smoothness is hardly given to the surface of the upper recording layer. Accordingly, although the above-mentioned method (i.e., method of using magnetic particles of the specific length in the longitudinal direction in each of the upper and lower recording layers) can improve the frequency characteristics of the resulting medium to a certain level, the frequency characteristics deteriorate in the high-frequency band in the case of an audio tape, and the S/N ratio decreases because of lowering of sensitivity and increase of noises in the case of a video tape."}
{"text": "The present invention relates to a method of controlling the position of cigarette filter paper.\nAs is known, in the manufacture of filter-tipped cigarettes, following formation of the actual cigarette consisting of normally white paper wrapped about the tobacco, a filter is positioned by way of an extension of the cigarette, and is joined to the cigarette by applying a connecting strip of normally yellow paper (known as \"tipping\" or \"filter paper\"). More specifically, the filter paper presents adhesive portions, and is wound about the filter material so as to slightly overlap the cigarette paper and so secure the filter to the cigarette.\nWhen applying the filter paper, however, this is sometimes positioned incorrectly in relation to the actual cigarette portion in the sense that the joining line between the cigarette and the filter is either crooked, or higher or lower than it should be.\nAs a result, a check must be made to control and indicate any errors in the position of the filter paper in relation to the cigarette."}
{"text": "1. Field of the Invention\nThis invention relates generally to a system or method of promoting tissue growth and more specifically to a multi-layer wound dressing having a tissue growth medium for enhancing the growth of tissue when exposed to reduced pressure.\n2. Description of Related Art\nVacuum induced healing of open wounds has recently been popularized by Kinetic Concepts, Inc. of San Antonio, Tex., by its commercially available V.A.C.® product line. The vacuum induced healing process has been described in commonly assigned U.S. Pat. No. 4,969,880 issued on Nov. 13, 1990 to Zamierowski, as well as its continuations and continuations in part, U.S. Pat. No. 5,100,396, issued on Mar. 31, 1992, U.S. Pat. No. 5,261,893, issued Nov. 16, 1993, and U.S. Pat. No. 5,527,293, issued Jun. 18, 1996, the disclosures of which are incorporated herein by this reference. Further improvements and modifications of the vacuum induced healing process are also described in U.S. Pat. No. 6,071,267, issued on Jun. 6, 2000 to Zamierowski and U.S. Pat. Nos. 5,636,643 and 5,645,081 issued to Argenta et al. on Jun. 10, 1997 and Jul. 8, 1997 respectively, the disclosures of which are incorporated by reference as though fully set forth herein.\nSubstantial work has also been performed relating to the creation of bioabsorbable and includable, cell growth enhancing matrices, lattices, or scaffolds. Exemplary U.S. patents known to applicant include Kemp et al. U.S. Pat. No. 5,256,418 issued Oct. 26, 1993; Chatelier et al. U.S. Pat. No. 5,449,383 issued Sep. 12, 1995; Bennett et al. 5,578,662 issued Nov. 26, 1996; and two patents issued to Yasukawa et al. U.S. Pat. No. 5,629,186 issued May 13, 1997 and U.S. Pat. No. 5,780,281 issued Jul. 14, 1998, both from a common parent application; the disclosures of which are incorporated by reference herein.\nAs is well known to those of ordinary skill in the art, closure of surface wounds involves the inward migration of epithelial and subcutaneous tissue adjacent the wound. This migration is ordinarily assisted through the inflammatory process, whereby blood flow is increased and various functional cell types are activated. Through the inflammatory process, blood flow through damaged or broken vessels is stopped by capillary level occlusion; thereafter, cleanup and rebuilding operations may begin. Unfortunately, this process is hampered when a wound is large or has become infected. In such wounds, a zone of stasis (i.e. an area in which localized swelling of tissue restricts the flow of blood to the tissues) forms near the surface of the wound.\nWithout sufficient blood flow, the epithelial and subcutaneous tissues surrounding the wound not only receive diminished oxygen and nutrients, but also are less able to successfully fight bacterial infection and thus are less able to naturally close the wound. Until the advent of vacuum induced therapy, such difficult wounds were addressed only through the use of sutures or staples. Although still widely practiced and often effective, such mechanical closure techniques suffer a major disadvantage in that they produce tension on the skin tissue adjacent the wound. In particular, the tensile force required in order to achieve closure using sutures or staples may cause very high localized stresses at the suture or staple insertion point. These stresses commonly result in the rupture of the tissue at the insertion points, which can eventually cause wound dehiscence and additional tissue loss.\nAdditionally, some wounds harden and inflame to such a degree due to infection that closure by stapling or suturing is not feasible. Wounds not reparable by suturing or stapling generally require prolonged hospitalization, with its attendant high cost, and major surgical procedures, such as grafts of surrounding tissues. Examples of wounds not readily treatable with staples or suturing include large, deep, open wounds; decubitus ulcers; ulcers resulting from chronic osteomyelitis; and partial thickness burns that subsequently develop into full thickness burns.\nAs a result of these and other shortcomings of mechanical closure devices, methods and apparatus for draining wounds by applying continuous negative pressures have been developed. When applied over a sufficient area of the wound, such negative pressures have been found to promote the migration toward the wound of epithelial and subcutaneous tissues. In practice, the application to a wound of negative gauge pressure, commercialized by Assignee or its parent under the designation “Vacuum Assisted Closure” (or “V.A.C.®”) therapy, typically involves the mechanical-like contraction of the wound with simultaneous removal of excess fluid. In this manner, V.A.C.® therapy augments the body's natural inflammatory process while alleviating many of the known intrinsic side effects, such as the production of edema caused by increased blood flow absent the necessary vascular structure for proper venous return.\nWhile V.A.C.® therapy has been highly successful in the promotion of wound closure, healing many wounds previously thought largely untreatable, some difficulty remains. Because the very nature of V.A.C.® therapy dictates an atmospherically sealed wound site, the therapy must often be performed to the exclusion of other beneficial, and therefore desirable, wound treatment modalities. One of these hitherto excluded modalities is the encouragement of cell growth by the provision of an in situ cell growth-enhancing matrix.\nAdditional difficulty remains in the frequent changing of the wound dressing. As the wound closes, binding of cellular tissue to the wound dressing may occur. Use of traditional V.A.C.® therapy necessitates regular changing of the dressing. Reckless dressing changes can result in some tissue damage at the wound site if cellular tissue has grown excessively into the dressing."}
{"text": "Articulated keyboard support mechanisms are typically mounted on and associated with work surfaces and desktops to support a computer keyboard. Such mechanisms are additionally often provided with an associated support pad for a user's wrist or palm. These support pads are conventionally fixed height foam or gel pads that lie along the front edge of the keyboard support on the upper surface thereof. While keyboarding, a user may rest his wrists or palms on the support pad to improve ergonomics. Different users require different heights in order to adequately support their hands. Moreover, recent developments in ergonomics suggest that a user should occasionally alter the position of their hands in an effort to reduce repetitive stresses. A fixed height support pad does not provide the desirable height adjustment relative to the keyboard. Moreover, keyboards may have different heights. Thus, one height of a support adjacent one keyboard may be inadequate adjacent a different keyboard. In fact, it may be desirable to remove the support from the keyboard support mechanism.\nNumerous keyboard support mechanisms have been developed for supporting a keyboard adjacent a worksurface. Some keyboard support mechanisms include a wrist or palm support fixed to the tray supporting the keyboard. In known arrangements, the height of the wrist support is adjustable to account for the differences in keyboard construction and users. For example, in some keyboard support mechanisms, the wrist support cooperates with two separate rotatable threaded members, which are rotatably secured to the keyboard support tray and individually manually rotated for adjusting the height of respective ends of the wrist support. Users of these types of mechanism experience difficulties in maintaining the support pad in a level orientation due to the difficulty in synchronizing movement of the two threaded members. Moreover, the threads typically require a plurality of time-consuming revolutions to adjust the height of the wrist support. In other more structurally complex mechanisms, the wrist support cooperates with a single lateral slide, which is slidably secured to the keyboard support tray and laterally slides to adjust the height of the wrist support. Examples of prior keyboard support mechanisms are found in: U.S. Pat. Nos. 4,545,554; 5,219,136; 5,375,800; 5,421,543; 5,443,237; 5,507,458; and 5,836,560.\nIt is an object of the invention to provide an improved keyboard support arrangement that accommodates both a keyboard support tray and a wrist/palm support pad while greatly improving upon both the structure and the convenience and flexibility of use in comparison to known structures.\nMore specifically, it is an object of the invention to provide a height adjustable wrist/palm support pad arrangement which is supported by a pair of rotatable camming hubs which are linked for synchronous rotational displacement and one of which has a manually-engagable actuator part to enable simple height adjustment by an operator.\nMore specifically, the present invention relates to a keyboard support arrangement which includes both a main keyboard support tray and a height adjustable wrist/palm support pad removably secured to the main keyboard support tray so as to provide increased convenience and flexibility of use with respect to the support pad. In the improved support pad arrangement, the support pad is easily selectively separable from or mountable on the keyboard support tray. A connecting assembly is provided on the pad for removably securing the support pad to the tray.\nIn accordance with a preferred embodiment of the invention, the keyboard support arrangement includes a keyboard support tray and a wrist/palm support pad adapted to support a user's wrist or palm thereon. A connecting and height adjusting structure detachably fixes the support pad on the tray upper surface and positions the support pad at one of a plurality of different heights. In one position of the connecting and height adjusting structure, the support pad is removable from the tray. In a second position, the support pad is secured to the tray at a first height. In a third position, the support pad is secured to the tray at a second height.\nFurther in accordance with the preferred embodiment of the invention, the connecting and height adjusting structure has two interconnected hubs each including radial flanges which are rotatably received in channels in the tray to rotatably fix the hubs to the tray. Once the support pad is attached to the tray, the connecting and height adjusting structure adjusts the height of the support pad relative to the tray by rotation of the hub causing opposed cam elements to react with and effect height adjustment of the support pad.\nAnother object of the invention is to provide a single actuator for both releasing the support pad from the keyboard support tray and adjusting the height of the support pad relative to the tray. More specifically, the rotatable hubs are linked together, and one of the hubs has a manually-engagable actuator part to effect synchronous rotation of the linked hubs to adjust the height of the support pad or release the support pad from the keyboard support tray.\nOther objects and purposes of the invention will be apparent upon reading the following specification and inspecting the accompanying drawings."}
{"text": "In the conventional two-stroke internal combustion engine, the inner surface of the cylinder is furnished with an exhaust port and a set of scavenging ports, which are controlled with the reciprocating piston to open and close. The upper edge of the exhaust port is above the upper edge of the scavenging ports so as to have the exhaust port opened first. After the power stroke of the aforesaid two-stroke internal combustion engine being completed, the waste gas in the cylinder is to be exhausted by means of gas mixture entered into the cylinder via the scavenging ports. In the two-stroke internal combustion engine that uses a carburetor, the air and fuel are pre-mixed into a homogeneous gas mixture, which is then conveyed into the cylinder, and which is also used for clearing the waste gas in the cylinder simultaneously.\nFor instance, a Japan patent application No. 51-39112 teaches a two-stroke internal combustion engine, in which a valve is furnished for opening or closing the exhaust port so as to have the height of the upper edge of the exhaust port become variable; the upper edge of the exhaust port is almost positioned above the upper edge of the scavenging ports.\nIn the aforesaid conventional two-stroke internal combustion engine, the exhaust port is opened earlier than the scavenging ports, but is closed later than the scavenging ports; therefore, a given portion of the gas mixture is to be blown off from the exhaust port upon the mixture entering into the cylinder via the scavenging ports; in other words, a considerable quantity of fuel is to be blown into the exhausting passage without being burned, that blow-off condition not only increases the fuel consumption, but also increases the HC quantity contained in the exhausted gas; that blow-off condition is particularly serious during the engine running in idle or low speed condition. That is the reason why the small vehicle with a two-stroke engine exhausts more HC in the waste gas during being stopped and the idle running period.\nIt is well known that the content of HC in the waste gas exhausted can be reduced by lowering the upper edge of the exhaust port to about the same level of the upper edge of the scavenging ports so as to reduce the blow-off volume; in that case, the scavenging efficiency of the engine will be reduced, and, as a result, the power output under normal using condition will be reduced; therefore, it is deemed impractical by lowering the upper edge of the exhaust port to prevent from leaking. Moreover, the height of the upper edge of the exhaust port may be changed by furnishing a valve, which can be opened at high speed running condition so as to make the upper edge of the exhaust port higher for increasing the power output; however, when the valve is under closed condition, the upper edge of the exhaust port is still at an ordinary height, being unable to prevent the blow-off condition from taking place.\nThere is an internal combustion engine that has no carburetor; the fuel is to be sprayed directly into the cylinder. During the scavenging phase, the waste gas can be exhausted through the exhaust port, and a portion of fresh air can enter into the cylinder via the scavenging port without losing fuel; however, the atomizing efficiency of the aforesaid engine is much worse than the engine with carburetor, such as the fuel not being sprayed evenly.\nThe atomized fuel drops have to be gasified before being burned. But in high speed running condition, the atomized fuel drops can be burned completely within a very short combustion period. As a result of the combustion speed being limited, the power can not increased; during the engine running at a high speed, a lot of atomized fuel drops would be exhausted out of the exhaust port before being burned thoroughly; therefore, a heavy fuel consumption and air pollution will be resulted."}
{"text": "1. Field of the Invention\nThis invention relates generally to sports towels, and more particularly, to a sports towel that can be used therapeutically.\n2. Background Information\nIt is generally well known that performing muscle-stretching exercises has wide-ranging health benefits. For example, getting blood flowing to the muscles by stretching, or “warming up,” before participating in sports can prevent injury to the muscles as well as enhance athletic performance. Further, sore or tightened muscles often benefit from stretching exercises after activity. It is also generally known that the application of ice can be helpful in promoting healing and reducing the pain and inflammation that can sometimes accompany tired or injured muscles.\nOften, many healthcare clinicians such as chiropractors and physical therapists prescribe stretching exercises and ice treatments to their patients. These clinicians may teach their patients to stretch using stretch tubing or stretch bands during a therapeutic office visit, to add resistance to the movement of the muscles in order to warm up the muscles properly. Many times physical therapists will schedule a series of muscle therapy sessions in which the patient uses the stretching devices, and then instruct the patient to continue to use the devices at home, in the gym, and on the sports field. Many office visits also include the application of ice, and the patient is often instructed to continue use of ice while at home.\nAlthough the benefits of stretching and icing are well known, compliance by patients and athletes tends to be poor. Very few people actually carry a muscle stretching device or an ice pack when exercising. Generally, this is because the muscle stretching devices and ice packs that are currently on the market are inconvenient for the participant to carry. It some cases, the public image of carrying therapeutic devices is also a reason for an athlete to leave the beneficial therapeutic devices behind. These factors are true for any number of sporting activities, such as cycling, football, soccer, weight lifting, general conditioning, etc.\nOne example of a sport that can particularly benefit from stretching and icing is golf. Performing stretching exercises prior to playing golf is recommended by golf experts as a means to enhance golf performance by increasing flexibility and range of motion. Typically, however, golfers do not properly stretch prior to playing golf. Many golfers simply take a few practice swings before hitting their first golf shot. These individuals tend to not get “warmed up” until after playing a few holes of golf, and often report “tightness” in their golf swing. Further, while on the golf course they may often suffer injuries, such as muscle pulls or sore backs and shoulders that would benefit from the application of ice after the round of golf is completed.\nThere remains a need, therefore, for a device that promotes stretching and icing of the muscles for therapeutic patients and athletes and is convenient to carry. There is also a need for a device that healthcare clinicians and trainers can prescribe to their patients for stretching and icing that promotes a continued use of the device."}
{"text": "Recently, inkjet printers have been developed which use printheads manufactured by micro-electro mechanical systems (MEMS) techniques. Such printheads have arrays of microscopic ink ejector nozzles formed in a silicon chip using MEMS manufacturing techniques. The invention will be described with particular reference to silicon printhead chips for digital inkjet printers wherein the nozzles, chambers and actuators of the chip are formed using MEMS techniques. However, it will be appreciated that this is in no way restrictive and the invention may also be used in many other applications.\nSilicon printhead chips are well suited for use in pagewidth printers having stationary printheads. These printhead chips extend the width of a page instead of traversing back and forth across the page, thereby increasing printing speeds. The probability of a production defect in an eight inch long chip is much higher than a one inch chip. The high defect rate translates into relatively high production and operating costs.\nTo reduce the production and operating costs of pagewidth printers, the printhead may be made up of a series of separate printhead modules mounted adjacent one another, each module having its own printhead chip. To ensure that there are no gaps or overlaps in the printing produced by adjacent printhead modules it is necessary to accurately align the modules after they have been mounted to a support beam. Once aligned, the printing from each module precisely abuts the printing from adjacent modules.\nUnfortunately, the alignment of the printhead modules at ambient temperature will change when the support beam expands as it heats up to the temperature it maintains during operation."}
{"text": "Heated fluids are often delivered to substrates, e.g. moving web-like substrates, for a variety of purposes. For example, heated fluids may be impinged upon a substrate for purposes of bonding, annealing, drying, promoting a chemical reaction, and the like."}
{"text": "1. Field of the Invention\nThis invention relates in general to garage door operators and in particular to a novel garage door operator wherein the receiver can be energized by two or more transmitted codes which are stored in the receiver.\n2. Description of the Prior Art\nGarage door operators of the prior art used transmitters in which the code can be changed by various methods as, for example, by moving two position switches to change the code. Such systems have also used code changing switches in the receiver so that the receivers can be set to correspond to the selected transmitter code.\nIt has also been known to use fixed frequency transmitters and fixed frequency receivers such that if the transmitted frequency matches the receiver frequency the receiver will respond.\nIt is an object of the present invention to provide a plurality of transmitters wherein each transmitter has its own unique and permanent non-user changeable code and wherein the receiver can be placed into a program mode wherein it will receive and store two or more codes corresponding to two different transmitters. The number of codes which can be stored in transmitters can be extremely high as, for example, greater than one million codes. Thus, the invention makes it possible to eliminate the requirements for code selection switches in the transmitters.\nIn the present invention the decoder module in the receiver will be capable of learning several different transmitted codes which will eliminate code switches in the receiver and also provides for multiple transmitters for actuating the garage opener.\nThe communication link can be various system such as radio frequency, light, wires, etc.\nThe invention makes it very easy for the user to operate the system and more secured code systems are available due to the higher number of available codes.\nAn encoded signal will be utilized wherein a pulse and blank time comprises a sync time base and different link pulses such as 1, 2, 3 or 4 milliseconds can be selected so as to provide different codings. Each datum can be 1, 2, 3 or 4 times the length of the sync pulse. The timing is from the rising edge to rising edges of the pulse and with ten data bits the number of codes can be in excess of one million codes.\nIn the invention, each transmitter encoder will contain a chip which contains a unique code and the receives will be able to memorize two or more as, for example, five different transmitter codes. This eliminates the need to have coding switches in either the transmitter or receiver. This eliminates the requirement that the user set the code switches so they match since the code switches are eliminated.\nIn the invention, during an operate mode, a receiver code must match an already programmed code four times in order to operate the garage door. This match is referred to as a valid code. Each valid code can be separated by up to two error codes and still have the output indicated as accurate.\nIn the program mode a code must be received four times in a row in order to be permanently stored in the receiver. Any error code will reset the valid code counter.\nThe advantage of the coding scheme are:\n1. Higher peak power without exceeding the FCC rules which gives longer transmitter range.\n2. Eliminate code switches in the transmitter and receiver making it easier for a customer to install and operate his garage door operator.\n3. Customers having more than one transmitter will not have to match codes.\n4. More secure codes due to the higher number of combinations which are available. Other objects, features and advantages of the invention will be readily apparent from the following description of certain preferred embodiments thereof, taken in conjunction with the accompanying drawings although variations and modifications may be effected without departing from the spirit and scope of the novel concepts of the disclosure and in which:"}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor devices and more particularly to a voltage-controlled self turn-off power semiconductor device with insulated-gate structure.\n2. Description of the Related Art\nPower semiconductor devices having higher withstanding or breakdown voltage characteristic are increasingly in demand in the recent manufacture of industrial-use equipment. For example, in the field of motor-controllers for railroad vehicles and industrial-use inverters, high-voltage strain switching devices with the withstanding voltage of more than several thousands volts in magnitude are required. Conventionally, self turn-off thyristors such as gate turn-off (GTO) thyristors are known as the semiconductor devices that fulfill such demand.\nThe self turn-off thyristor is principally one of the current-controlled switch devices. The thyristor will latch up when it turns on, and has an advantage of attaining a decreased \"on-resistance\" thereof. Decrease in the on-resistance potentially decreases the on-voltage of the thyristor, causing its on-characteristic to be improved. On the other hand, this type of thyristor remains inherently low in the maximum cut-off current density. In particular, a certain type of GTO thyristor that turns off using an insulated gate or metal-oxide-semiconductor (MOS) structure, such as a MOS-controlled thyristor (MCT), suffers from the poor current cut-off ability. This means that such type of thyristor should have a limit in its high-voltage strain characteristic. The breakdown voltage limit is a serious bar to the wide applicability of the GTO thyristor as power controller devices.\nIn the recent past, there has been proposed and developed a bipolar-type semiconductor device that includes a bipolar transistor section, which is driven to turn on and off by an insulated gate structure. This type of device is generally known as the \"insulated-gate bipolar transistor (IGBT)\". With a presently available IGBT, a MOS transistor having an insulated gate electrode is added to the base of the transistor section, thereby to enable switching operations to be performed by potentially changing the insulated gate electrode. This means that the IGBT is principally a voltage-controlled switch device. In this respect, the IGBT may be more suitable than the GTO thyristor in applicability as an industrial-use high-voltage strain switch device. Unfortunately, while the IGBT is high in the maximum cut-off current density achievable, it suffers from a serious disadvantage that the on-resistance stays high. This is due to the technical limitation that the IGBT device must be so designed that a latch-up will not occur in an internal thyristor structure, which is constituted by an alternate lamination (PNPN) structure of P- and N-type layers inherent within the IGBT.\nThe above analysis teaches that the presently available GTO thyristors or insulated-gate GTO thyristors are low in the maximum cut-off current density, whereas IGBTs are high in the on-resistance, and that each of the devices is disadvantageous in that they fail to meet the essential requirements for the high demand self turn-off power semiconductor devices. To date, a self turn-off power semiconductor device has not been achieved yet which can attain both a decreased on-resistance and an increased maximum cut-off current density (enhanced withstanding voltage characteristic)."}
{"text": "This relates generally to strain gauges, and, more particularly, to strain gauges for monitoring strain in printed circuits.\nElectronic devices such as cellular telephones, computers, and other electronic devices contain integrated circuits and other electrical components. Components such as these may be mounted on printed circuits. During drop events and other situations in which an electronic device is subjected to an impact or other conditions leading to elevated stresses, solder joints may be weakened and other faults can develop in the circuitry mounted on a printed circuit. Unless care is taken, stress-induced damage to an electrical component or other circuitry in an electronic device may create reliability issues.\nTo help understand the way in which strain is distributed to the components in a printed circuit, strain tests may be performed on a test printed circuit board. External strain gauges may be attached to the upper and lower surfaces of the test board in the vicinity of integrated circuits or other components of interest. These strain gauges may be wired to test equipment that gathers strain data. Strain data may be gathered with the strain gauges while applying force to the printed circuit. By analyzing the strain data, printed circuit board designs and component layouts can be refined to enhance reliability.\nIt can be difficult to gather accurate strain data with this type of approach. Strain measurements from strain gauges mounted on the surface of the printed circuit board may be influenced by the way in which the strain gauges are adhered to the surface of the printed circuit board and other variables that are difficult to control. These strain gauges also do not gather strain data in the field to alert a user or others about the presence of excessive strain.\nIt would therefore be desirable to be able to provide improved strain gauge configurations for monitoring components on printed circuit boards."}
{"text": "(a) Field of the Invention\nThe present invention relates to a method for the production of alkylates by sulfuric acid alkylation of isoparaffins with olefins. This method is particularly well adapted for use in the petroleum refining industry, using isobutane as isoparaffin, and butylenes as olefin.\nThe invention also relates to a device for mixing and reacting at least two and preferably three liquid components. This device is particularly well adapted for carrying out the above method even though it can be used for carrying out many other methods.\n(b) Brief Description of the Prior Art\nAlkylates are main components of high-octane motor fuels. They are produced by alkylation of isoparaffins (mainly isobutane) by olefins (such as propylene, butylenes or amylenes) in the presence of sulfuric or hydrofluoric acid that serves as a process catalyst. The most widely known method for the production of alkylates in the petroleum refining industry consists of carrying out a sulfuric acid alkylation of isobutane by olefin.\nNumerous methods for carrying out sulfuric acid alkylation of isobutane by olefins are known. The method according to the invention distinguishes over most of these known method in that the reaction is carried out in a compact reactor which does not contain moving parts and in which jet mixing of the reagents is achieved.\nU.S. Pat. No. 3,544,652 issued on Dec. 1st, 1970 discloses a method for the alkylation of isoparaffin by olefins in the presence of sulfuric acid, where the olefin is reacted with an alkylating hydrocarbon-in-acid emulsion formed by thoroughly mixing isoparaffin with sulfuric acid before contact with the olefin. In this patent, the isoparaffin-to-olefin volume ratio is disclosed as being equal to about 12:1. The acid-to-hydrocarbons volume ratio is disclosed as being within the range of 2.5:1 to 15:1 but it is mainly maintained at about 6:1. The reaction is carried out adiabatically, mainly in a continuous manner, in a reactor called “alkylation contactor”, which is provided with a mixer that is devised for forming the isoparaffin-in-sulfuric acid emulsion and for thoroughly and homogeneously mixing the so-formed emulsion with the olefin at the points of delivery of the latter into the reactor.\nAs the liquid flows through the reactor, the temperature of the alkylating mixture rises continuously by 5 to 15° C., thereby reducing viscosity of the mixture and increasing its turbulence. The method is carried out at a temperature of 5 to 60° C. under a pressure sufficient for keeping the reagents in a liquid state (from 2 to 10 ATMs). Prior to being mixed with the isoparaffin, sulfuric acid at a concentration of 88 to 99% is cooled down to a temperature of about 4° C.\nThe emulsion preparation and the olefin injection and distribution inside the reaction area are not disclosed in detail in this U.S. patent.\nThe method disclosed in U.S. Pat. No. 3,544,652 is efficient but it requires a substantial amount of power for circulating the acid due to the very high acid-to-hydrocarbons ratio. It also requires a settling equipment of a very large size. Moreover, the method disclosed in this patent cannot guarantee a low consumption of sulfuric acid and a reasonably high quality of the final product.\nRussian patent No. 2,131,861 granted on Jul. 25, 1994 (corresponding to U.S. Pat. Nos. 5,443,799 and 5,777,189) discloses a method for sulfuric acid alkylation of isoparaffins by olefins and a device for carrying out this method. At the initial stage of the method disclosed in this patent, a thin isoparaffin-in-sulfuric acid emulsion is made by injecting isoparaffin into an acid medium through a set of nozzles. Then the emulsion is delivered into a reaction area where olefin is fed, through a number of points normal to the emulsion flow. In this method, the alkylation is carried out under isobaric and isothermal conditions.\nRussian patent No. 2,131,861 also discloses that the emulsion should preferably flow in the emulsion area at a rate of 0.2 to 2 m/s—and within the reaction area at a rate of 0.04 to 0.27 m/s. Depending on the selected flow rate, the contact between the reagents may last from a few to 60 seconds, thereby reducing to a minimum the possibility of not-wanted side reactions such as oligomerization of olefins and autoalkylation of isoparaffins. Tests have shown that this method permits to prepare a thin unstable emulsion. Separation of the reaction mixture into a hydrocarbon phase and an acid phase takes 5 to 8 seconds, thereby allowing reduction in the setting time.\nSince the method described in Russian patent No. 2,131,861 does not require rotary mixers, the equipment required for carrying it out is rather cheap and of easy control and maintenance.\nRussian patent No. 2,131,861 further discloses a device for carrying out the above method. This device comprises a tank for preparing the emulsion. A special appliance is provided for isoparaffin injection within the tank. Such an appliance essentially consists of a set of axially arranged nozzles. An appliance is also provided for sulfuric acid injection within the tank. The device also comprises a mixing chamber that is part of the tank, with an outlet throat, and a cylindrical reactor which is connected in line to the throat of the emulsion preparation tank. To provide olefin injection, the device comprises a perforated branch pipe extending along the axis of the reactor.\nThe method and device described in the above Russian patent No. 2,131,861 and its foreign counterparts have rather acceptable technical and economic parameters of operations, as proved by industrial tests. However, those parameters could be improved if use is made of a higher level of flow turbulence in the reaction zone. In practice, such a higher level of flow turbulence could be obtained if the flow rate is increased in the reaction zone and the mixing conditions of olefin and emulsion flows are improved by using a more efficient olefin feed unit instead of using a perforated branch pipe extending along the axis of the reactor.\nRussian patent No. 2,092,475 granted on Dec. 6, 1995 is the closest prior art known to the Applicant. It discloses a method for the production of alkylates in a tubular reactor, which comprises mixing sulfuric acid with isobutane previously cooled down to a temperature of not over −2° C.; mixing the obtained emulsion with olefins also previously cooled to a temperature of not over −2° C., in a plurality of stages; separating the sulfuric acid from the obtained reaction mass; and recycling it. This method requests that the sulfuric acid be mixed with the isobutane and the obtained emulsion be mixed with the olefin in an injector-type mixer, with an isobutane-to-sulfuric acid injection ratio of 3.3 to 5.2 and an isobutane-to-olefin volume ratio of 3000-5000:1. In this method, sulfuric acid is separated from the reaction mass in a hydrocyclone.\nRussian patent No. 2,092,475 also discloses a device for carrying out the above method, which consists of a reactor provided with three concatenated injection mixers. Each mixer is provided with an olefin injection appliance that distributes the feed in a helical fashion along the length of a device.\nWith the method and device disclosed in Russian patent No. 2,092,475, one may carry out sulfuric acid alkylation of isoparaffins by olefins in a compact reactor that has no moving parts. One may also obtain a high quality alkylate, as proved by industrial tests. However, in practice, an isobutane-to-olefin ratio of 3,000-5,000:1 can be obtained in the reaction area only with very high power consumption. Alkylate quality and specific sulfuric acid consumption could actually be improved by using a more efficient olefin feed unit than the perforated tube with helical openings as described in this patent. Besides, alkylate quality and specific sulfuric acid consumption could be improved by optimization of the alkylation process conditions, by proper selection of different size ratios of the reactor elements and by an improved process of separation of the reaction mixture and a device to carry out this process."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor light emission element using a nitride semiconductor, a semiconductor composite element and a process for producing a semiconductor light emission element.\n2. Description of the Related Art\nWide band gap semiconductors have been receiving attention as a semiconductor for a blue light emission element in recent years (S. Nakamura, T. Mukai and M. Senoh, Jpn.J. Appl. Phys., vol. 30, L1998 (1991); and S. Nakamura, M. Senoh, S. Nagahama, N. Iwata, T. Yamada, M. Matsushita, Y. Sugimoto and H. Kiyoku, Appl. Phys. Lett., vol. 69, p. 1477 (1996)). Among them, nitride semiconductors (nitride compound semiconductors) are receiving attention as the wide band gap semiconductor.\nHowever, the nitride semiconductors are often produced by a two-step growth, in which a low temperature buffer layer is firstly formed by using a sapphire substrate, and after increasing the temperature to about 1,000xc2x0 C., crystal growth is carried out at a high temperature of 1,000xc2x0 C.\nWith respect to silicon semiconductors, which constitute major portion of semiconductor elements, such as LSI, there has been an attempt that light is emitted from silicon to carry out transmission between chips and within a chip with light (K. D. Hirschman, L. Tsybeskov, S. P. Duttagupta and P. M. Fauchent, Nature, vol. 384, p. 338 (1996) and S. Fujita and N. Sugiyama, Appl. Phys. Lett., vol. 74, p. 308 (1999)).\nHowever, silicon cannot emit light as it is because it is an indirect transition semiconductor. Therefore, direct transition III-V compound semiconductors (nitride semiconductors), which emit light, are formed on silicon, but it is the current situation that there is a problem in that crystals of high quality cannot be grown. It is considered that this is because there are excess bonding linkages left severed (dangling bonds) at the junction interface of the two different semiconductors having different lattice constants and they form interface levels. Particularly, in the case where there is a large difference in thermal expansion coefficient between the two semiconductors, it is considered that large stress is formed after cooling because a high temperature is applied upon forming the nitride semiconductor, whereby not only the formation of junction defects is accelerated, but also structural defects, such as cracking, crazing and interfacial exfoliation, are caused.\nIt is the current situation of the semiconductor light emission element obtained by forming the direct transition III-V compound semiconductors (nitride semiconductor), which emit light, on silicon, that the characteristics of the junction interface are deteriorated by the formation of a large amount of interfacial level at the junction interface, and sufficient characteristics as a semiconductor light emission element cannot be brought out.\nSuch a semiconductor light emission element is also studied that a nitride semiconductor is formed on a dissimilar semiconductor substrate formed with inexpensive Si or GaAs, which is different from the sapphire substrate or the nitride semiconductor. However, two-step growth is carried out for growing the nitride semiconductor by using a buffer layer formed with GaN, AlN or ZnO, and therefore, the semiconductor substrate formed with the dissimilar semiconductor, such as Si and GaAs, functions only as a substrate to grow the nitride semiconductors but does not serve as a semiconductor light emission element. (Supratik Guha and Nestor A. Bojarczuk, Appl. Phys. Lett., vol. 72, p. 415 (1998), Chuong A. Tran, A. Osinski and R. F. Karlicek, Jr., Appl. Phys. Lett., vol. 75, p. 1494 (1999), and E. Feltin, S. Dalmasso, P. Mierry, B. Beaumont, H. Lahreche, A. Bouille, H. Haas, M. Leroux and P. Gibart, Jpn. J. Appl. Phys., vol. 40, L738 (2001))\nWhile a heterogeneous junction between nitride semiconductors has been realized by a quantum well structure, the heterogeneous junction to the dissimilar semiconductor, such as Si and GaAs, is difficult to be formed as noted in the foregoing. Furthermore, there are such problems that the nitride semiconductor crystals are formed by providing a buffer layer, such as GaN, AlN and ZnO, at the interface to the dissimilar semiconductor, and the semiconductor junction cannot be directly formed.\nAs described in the foregoing, in the semiconductor light emission element having a constitution containing a nitride semiconductor layer and a dissimilar semiconductor (semiconductor substrate) different therefrom, the nitride semiconductor layer and the dissimilar semiconductor (semiconductor substrate) are formed through the buffer layer, such as GaN, AlN and ZnO, having insulating property, and the dissimilar semiconductor does not contribute to injection of charge (electrons or positive holes), whereby a large number of layers are required, and problems in cost are caused. Furthermore, in the case where the dissimilar semiconductor (semiconductor substrate) is shared by the other elements, and driving circuits and wiring from the other elements are formed on the semiconductor substrate in parallel with the semiconductor light emission element, the elements are not formed directly on the semiconductor substrate, and the semiconductor substrate does not electrically function by itself. As a result, there are problems of complex layer structure and high production cost, and it is the current situation that a composite element that has plural elements having different functions is difficult to be formed on the same semiconductor substrate.\nThe invention is to solve the foregoing problems associated with the conventional art to provide a simple and low cost semiconductor light emission element exerting high performance and a process for producing the same.\nThe invention is also to provide such a semiconductor light emission element capable of forming a semiconductor composite element that has plural semiconductor elements having different functions on the same dissimilar semiconductor (semiconductor substrate) and has high performance and a simple structure, and to provide the semiconductor composite element.\nThe invention provides, as one aspect, a semiconductor light emission element containing a nitride semiconductor layer at least containing one or more elements selected from Group IIIA elements and one or more elements selected from Group VA elements, a semiconductor having a polarity different from the nitride semiconductor layer and being dissimilar to the nitride semiconductor layer, and a light emission layer provided between the dissimilar semiconductor and the nitride semiconductor, wherein electrons or positive holes are injected from semiconductors of the dissimilar semiconductor and the nitride semiconductor layer to the light emission layer to carry out light emission."}
{"text": "1. Technical Field\nThe present disclosure relates to face recognition technologies and, particularly, to a face recognition system and a face recognition method used for taking a group photo.\n2. Description of Related Art\nIt is difficult to take a good group photo, especially when the group is too big to make sure that all people's faces in the group are completely taken. Therefore, it is desirable to provide a face recognition system and a face recognition method, which can overcome the above-mentioned shortcomings."}
{"text": "It is known to use ultrasound to image anatomical structures by transmission and reception of ultrasound waves from a transducer. Ultrasound imaging may often be applied to image the fetus in the womb.\nThree-dimensional (3D) ultrasound images may be obtained by using software to combine ultrasound data that has been taken at different positions or angles to obtain volumetric ultrasound data, and to render an image from the volumetric ultrasound data using methods such as simple surface shading or direct volume rendering. The rendered image may be represented as a 2D rendered image data set, which comprises appropriate pixel color values for display on a 2D surface such as a 2D screen. Such rendered images may nonetheless be referred to as 3D images because they provide an image that the viewer perceives as being representative of a structure in three dimensions, for example a solid structure.\nIn four-dimensional (4D) ultrasound imaging systems, a series of 3D images obtained at different times is dynamically rendered to produce a moving 3D image, for example a 3D ultrasound movie.\nIn recent years, 3D and 4D ultrasound images have been made more realistic through the use of advanced lighting techniques (referred to, for example, as global illumination, gradient free lighting, subsurface scattering or photon mapping) that simulate illumination with a more physically accurate model than was previously used.\nIn global illumination, a lighting model may be used that includes both direct illumination by light coming directly from a light source and indirect illumination, for example illumination by light that has been scattered from another surface.\nGlobal illumination in ultrasound has gone from being seen as a novelty to being, in some cases, a required feature of ultrasound systems. In obstetrics, global illumination may be used to provide high quality images of the fetal face. Global illumination may also be of interest to doctors in other specialties, For example, global illumination may be used to render images from abdominal scans, so as to produce images that may in some cases be easier to read and/or contain greater detail than those produced with simpler lighting models. It is expected that the recognition and importance of global illumination may continue to increase.\nVolume rendering may be computationally costly. FIG. 1 shows a cubic grid of voxels for which width=height=depth=N. The complexity of the volume rendering for this cube may be of the order of N3 (all the voxels may have to be taken into account).\nVolume rendering for which global illumination is used may be even more computationally costly than volume rendering without global illumination. A naïve method for simple lighting configurations may include of the order of N4 operations.\nThere are global illumination methods that transform the global illumination process into two passes. Using such methods, the complexity of rendering may be considered to be reduced to be on the order of 2*N3. One such method is photon mapping. Other methods are also available that perform two passes, which may be referred to as traversals.\nA first pass may create a light volume, and a second pass may use the light volume to render an image for display. A two-pass system may allow a single light volume to be used for multiple rendering passes at different viewing angles, thus avoiding the need to recalculate the light volume for each render.\nA volumetric image data set comprises an array of voxels, each with an associated intensity. A volumetric image data set may be obtained directly from a medical imaging scan or through further processes such as reconstruction and filtering.\nThe volumetric image data set may be representative of all or part of the subject of the scan. The intensities of the voxels may correspond to physical properties of the subject of the scan, for example tissue types. The intensities may be mapped on to opacity and color values, for example by using a transfer function.\nThe volumetric image data set contains a representation of materials, surfaces and so on that occur in the subject. In the discussion below, processes may be referred to as if they occurred in a physical space (for example, light reflecting from a surface). However, in the case of illumination, we are usually describing virtual (simulated) processes occurring as numerical operations on a volumetric image data set. Similarly, when we discuss the volumetric image data set as if it were a physical space having a physical extent, we are generally referring to the coordinate space that is represented by the voxels of the image volume.\nIn order to render an image from the volumetric image data set, the position of at least one light source may be determined with reference to the volumetric image data set. The position of a viewpoint (which may be referred to as a camera) may also be determined. The image will be rendered as if viewed from the viewpoint.\nA first pass may comprise a traversal from the light source into the volumetric image data set, in which virtual light is cast into the volumetric image data set. A global illumination lighting model may be used. The irradiance due to the light source may be determined at each of a large array of points in the volumetric image data set using absorptive properties assigned to the voxels in dependence on the voxel intensities.\nThe irradiance values at the array of points may be stored as a light volume. The light volume may be stored in memory. The light volume may be independent of the viewpoint.\nThe irradiance of each voxel in the light volume is determined by the number of photons that passes through the voxel. For better numerical performance, photons are often treated in a statistical way, rather than in a discrete way. This means that light from the light source may act as a beam of light sometimes, but like a discrete photon at other times.\nAbsorption may be more efficient to describe statistically, usually as an integral formed by the Beer-Lambert law. Absorption as defined by the Beer-Lambert law is governed by a medium-specific coefficient called the attenuation coefficient. The attenuation coefficient has a range of [0,∞] and may be used synonymously with opacity. Opacity is often defined on a [0,1] range, indicating what fraction of the light is absorbed.\nReflection and scattering and refraction may be seen as events that can happen along the light path and as such may be part of the determination of each voxel's irradiance.\nA second pass may comprise a traversal through the light volume from the camera, using the light volume to provide global lighting information. Rays may be cast from the camera (for example, one ray for each pixel of the resulting rendered image), and irradiances from points along each ray may be integrated to provide pixel color values for a final rendered image."}
{"text": "The present invention concerns post-treated titanium dioxide (TiO.sub.2) having pigmentary morphological and granulometric characteristics, as well as a process for producing said product.\nOne object of this invention is to provide a TiO.sub.2 particle with pigmentary morphological and granulometric characteristics and endowed with a new property: i.e.; to have a chemically reactive coating.\nStill another object of this invention is that of conferring to such a coating a strong adherence to the TiO.sub.2 particles.\nYet another object of this invention is that of imparting to such a coating the property of being uniform and compact.\nA further object of the invention is that of conferring to the TiO.sub.2 particles thus coated, a specific surface controllable at will.\nAnother object still, is that of providing a process for obtaining the above said chemically reactive coating.\nStill another object of the invention is that of preparing a pigment superfically coated with organic molecules which impart to it an affinity for the organic substances of the vehicles, for example, paint vehicles, starting from the TiO.sub.2 with a chemically reactive coating.\nAll these objects and others still are achieved by the new product object of this invention, which consists of TiO.sub.2 particles with pigmentary morphological and granulometric characteristics and is characterized in that said product has a chemically reactive coating consisting of a mixture of oxides and oxychlorides selected from the group consisting of Al.sub.2 O.sub.3 and AlOCl; SiO.sub.2 and SiOCl.sub.2 ; and ZrO.sub.2 and ZrOCl.sub.2.\nThe chemical reactivity of these new products is due to the presence of the oxychlorides.\nAs a result of these reactive groups, one may, for example, fix onto the particles substances containing amine, carboxylic, or hydroxylic groups, thereby obtaining pigments with an affinity for organic vehicles, for example, vehicles for paints based on alkyd, polyester, melaminic, acrylic and phenolic resins.\nIn general, the atomic ratio between Cl and Al, Si or Zr, present in these new products, is between 0.10 and 0.90.\nStill another object of this invention is a process for preparing TiO.sub.2 particles with pigmentary morphological and granulometric characteristics having the above mentioned chemically reactive coating, this process being characterized in that the surface of the particles is activated in a moving bed by thermal treatment with an anhydrous gas at temperatures comprised between 400.degree. and 600.degree. C., and that the particles are then reacted in a moving bed with an inorganic chloride chosen out of a group consisting of: AlCl.sub.3, SiCl.sub.4 and ZrCl.sub.4, in the presence of a carrier gas, at temperatures between 350.degree. an 600.degree. C."}
{"text": "A portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent document or the patent disclosure, as it appears in the Patent and Trademark Office patent files or records, but otherwise reserves all copyrights whatsoever.\nIn a binary search, the conventional method used is to determine the midpoint of a sorted list of items of data or a subset of a list by calculating the sum of the extreme pointers and then dividing by two. In a conventional binary search half of the remaining list is eliminated from the search by each comparison. When the searched item is not found, the pointers cross to indicate where that searched item would be inserted to be in sorted order. Given N items in a list, N/2 items on average must be physically moved to make room for a new item to be inserted in sorted order."}
{"text": "Often it is desired to apply information to large or irregular items. For example, it is often desirable to apply postage indicia and other information, such as a mailing address, to parcels or to envelopes which have been filled with items to be mailed. However, parcels are typically too large and vary in size and shape too much from parcel to parcel to enable adapting a typical postage meter mechanism or other printer for printing postage indicia on parcels. Envelopes which have been filled with items to be mailed present items of varying thickness and which tend to result in non-flat surfaces making it difficult to adapt a typical postage meter mechanism or other printer for printing postage indicia on such envelopes.\nTraditional postage meters have been unable to print postage indicia directly upon parcels. Instead, adhesive labels have been used, wherein a label is inserted into a postage meter for printing of a postage indicia thereon. The adhesive label may then be applied to a surface of the appropriate parcel.\nMany traditional postage meters have been adapted to apply postage indicia to an envelope which has been filed with items to be mailed. For example, such postage meters have included a slot capable of accepting envelopes of varying thickness up to 5/16ths of an inch. A lower platen in the envelope printing area may be adapted to lift the envelope and apply force in the direction of an impact print mechanism. An impact print mechanism, which typically comprises an inked die shaped to print the desired postage indicia, is moved rapidly toward the face of the envelope (moved in the direction of the platen supporting the envelope) to impress a postage indicia image thereon. Such mechanisms have been successful in applying relatively simple postage indicia and/or other information to items of varying thickness and slightly non-flat surfaces. However, such mechanisms have not been successful with respect more complex printing operations, such as information based indicia (IBI) having machine readable barcode therein.\nAttempts have been made to adapt printers capable of printing more complex postage indicia, such as the aforementioned IBI, for printing upon envelopes which have been filled with items to be mailed. Such attempts have included ink jet printing mechanisms using a lifting platen and a print head lowering assembly to accommodate envelopes of varying thickness up to ⅜ths of an inch. However, such mechanisms are complicated, requiring both a mechanism to raise the lower platen and to lower the print head. Moreover, such mechanisms do not address the problem of the surfaces of the envelope which has been filled with items to be mailed not being flat, particularly at or near the edges of the envelope. Such mechanisms tend to provide acceptable printing of complex postage indicia only in areas of the envelope surface where it is relatively flat, and generally result in a higher than traditional number of unusable postage indicia (e.g., a barcode is deformed or imperfect to the point that the information therein cannot be read by a machine).\nAs can be appreciated from the above, postage meters or other printers generally available for use in printing postage indicia have heretofore been unable to print postage indicia or other information directly upon parcels. Postage meters or other printers available for use in printing postage indicia have been limited with respect to the range of envelope thicknesses which can be accommodated and have met with only limited success with respect to printing complex postage indicia. Moreover, such postage meters or other printers have been complicated and expensive, despite their limitations with respect to the sizes and shapes of items that may be printed upon."}
{"text": "Wind power is considered one of the cleanest, most environmentally friendly energy sources presently available, and wind turbines have gained increased attention in this regard. A modern wind turbine typically includes a tower, a generator, a gearbox, a nacelle, and a rotor including one or more rotor blades. The rotor blades capture kinetic energy from wind using known foil principles and transmit the kinetic energy through rotational energy to turn a shaft coupling the rotor blades to a gearbox, or if a gearbox is not used, directly to the generator. The generator then converts the mechanical energy to electrical energy that may be deployed to a utility grid.\nDuring operation, the direction of the wind which powers the wind turbine may change. The wind turbine may thus adjust the nacelle through, for example, a yaw adjustment about a longitudinal axis of the tower to maintain alignment with the wind direction. In addition, wind turbines typically use pitch control to maintain rated power in high winds. As wind speeds increase, the blades are “pitched to feather” (i.e. the pitch angle of the rotor blades is increased) to reduce lift and thrust and maintain torque and power of the wind turbine. When the rotor blades pitch excessively in high winds, aerodynamic blade noise can be increased as well, as a result of thickened boundary layers or separated flow on the pressure side surfaces of the rotor blades.\nAccordingly, it would be advantageous to limit excessive pitching and avoid an undesired associated noise increase. As such, the present disclosure is directed to systems and methods which actively yaw the nacelle of the wind turbine away from the nominal wind direction at high wind speeds. With the turbine yawed out of the wind, there is less available power and the blades will not be forced to pitch as excessively to maintain rated power."}
{"text": "The present invention relates to crawling (or traversing) of hyperlinked documents. More specifically, the invention relates to techniques for the distributed crawling of hyperlinked documents that can perform rate limiting of hosts and adapt to actual retrieval times of the hosts.\nThe World Wide Web (or “Web”) contains a vast amount of information in the form of hyperlinked documents (e.g., web pages). One of the reasons for the virtually explosive growth in the number of hyperlinked documents on the Web is that just about anyone can upload hyperlinked documents, which can include links to other hyperlinked documents. Although there is no doubt that there is a vast amount of useful information on the Web, the unstructured nature of the Web can make it difficult to find the information that is desired.\nSearch engines allow users to enter queries (e.g., key words) that describe the information users are seeking. The search engines then scan the Web for hyperlinked documents that best satisfy the query. With literally millions of hyperlinked documents on the Web, web crawlers are typically utilized to scan, index and store information regarding hyperlinked documents on the Web so that the search engines can execute queries more efficiently.\nAs the size of the Web continues to increase, it becomes increasingly more desirable to have innovative techniques for efficiently crawling the Web. Additionally, it would be beneficial to have web crawling techniques that are efficient yet do not impose unnecessary burdens on hosts on the Web."}
{"text": "Bone grafts are commonly used to treat defects in the skeletal system caused by injury or disease. Skeletal defects often require bone grafts to maintain space and provide a matrix for healing. A graft should provide or facilitate the various mechanisms of bone healing including osteoconduction, osteoinduction, and osteogenesis. Osteoconduction is the ability of the graft to act as a matrix or scaffold to support bone formation. Osteoinduction is a result of bone growth factors that stimulate differentiation of local cells to become bone forming cells, i.e. osteoblasts. Bone morphogenic proteins (BMP's) that are naturally occurring in bone, or that may be produced by recombinant gene technologies, are responsible for osteoinduction. Osteogenesis refers to the formation of bone, and may also be used to reference the ability of cells, to form bone. Bone forming cells may either be resident at the graft site or transplanted to the site by autogenous bone, bone marrow aspirate and/or cell implantation. Considering these requirements to form bone, a need exists for a reproducible and cost-effective process of making a bone graft having improved osteoconductive and osteoinductive properties."}
{"text": "1. Field of the Invention\nThis invention relates to enzyme inhibitors, and more particularly, to novel substituted biaryl oxobutyric acid compounds or derivatives thereof useful for inhibiting matrix metalloproteases.\n2. Description of the Related Art\nThe matrix metalloproteases (a.k.a. matrix metalloendo-proteinases or MMPs) are a family of zinc endoproteinases which include, but are not limited to, interstitial collagenase (a.k.a.. MMP-1), stromelysin (a.k.a.. proteoglycanase, transin, or MMP-3), gelatinase A (a.k.a.. 72 kDa-gelatinase or MMP-2) and gelatinase B (a.k.a. 95 kDa-gelatinase or MMP-9). These MMPs are secreted by a variety of cells including fibroblasts and chondrocytes, along with natural proteinaceous inhibitors known as TIMPs (Tissue Inhibitor of MetalloProteinase).\nAll of these MMPs are capable of destroying a variety of connective tissue components of articular cartilage or basement membranes. Each MMP is secreted as an inactive proenzyme which must be cleaved in a subsequent step before it is able to exert its own proteolytic activity. In addition to the matrix destroying effect, certain of these MMPs such as MMP-3 have been implemented as the in vivo activator for other MMPs such as MMP-1 and MMP-9 (Ito, et al., Arch Biochem Biophys. 267, 211 (1988); Ogata, et al., J. Biol. Chem. 267, 3581 (1992)). Thus, a cascade of proteolytic activity can be initiated by an excess of MMP-3. It follows that specific MMP-3 inhibitors should limit the activity of other MMPs that are not directly inhibited by such inhibitors.\nIt has also been reported that MMP-3 can cleave and thereby inactivate the endogenous inhibitors of other proteinases such as elastase (Winyard, et al., FEBS Letts. 279, 1, 91 (1991)). Inhibitors of MMP-3 could thus influence the activity of other destructive proteinases by modifying the level of their endogenous inhibitors.\nA number of diseases are thought to be mediated by excess or undesired matrix-destroying metalloprotease activity or by an imbalance in the ratio of the MMPs to the TIMPs. These include: a) osteoarthritis (Woessner, et al., J. Biol.Chem., 259(6), 3633 (1984); Phadke, et al., J. Rheumatol. 10, 852 1983)), b) rheumatoid arthritis (Mullins, et al., Biochim. Biophys. Acta 695, 117 (1983)); Woolley, et al., Arthritis Rheum. 20, 1231 (1977); Gravallese, et al., Arthritis Rheum. 34, 1076 (1991)), c) septic arthritis (Williams, et al., Arthritis Rheum. 33, 533 (1990)), d) tumor metastasis (Reich, et al., Cancer Res., 48, 3307 (1988), and Matrisian, et al., Proc. Nat\"\"l. Acad. Sci., USA 83, 9413, (1986)), e) periodontal diseases (Overall, et al., J. Periodontal Res. 22, 81 (1987)), f) corneal ulceration (Burns, et al., Invest. Opthalmol. Vis. Sci. 30, 1569 (1989)), g) proteinuria (Baricos, et al., Biochem. J. 254, 609 (1988)), h) coronary thrombosis from atherosclerotic plaque rupture (Henney, et al., Proc. Nat\"\"l. Acad. Sci., USA 88, 8154 (1991)), I) aneurysmal aortic disease (Vine, et al., Clin. Sci. 81, 233 (1991)), j) birth control (Woessner, et al., Steroids 54, 491 (1989)), k) dystrophobic epidermolysis bullosa (Kronberger, et al., J. Invest. Dermatol. 79, 208 (1982)), and I) degenerative cartilage loss following traumatic joint injury, m) conditions leading to inflammatory responses, osteopenias mediated by MMP activity, n) tempero mandibular joint disease, o) demyelating diseases of the nervous system (Chantry, et al., J. Neurochem. 50, 688 (1988)).\nThe need for new therapies is especially important in the case of arthritic diseases. The primary disabling effect of osteoarthritis (OA), rheumatoid arthritis (RA) and septic arthritis is the progressive loss of articular cartilage and thereby normal joint function. No marketed pharmaceutical agent is able to prevent or slow this cartilage loss, although nonsteroidal anti-inflamnmatory drugs (NSAIDs) have been given to control pain and swelling. The end result of these diseases is total loss of joint function which is only treatable by joint replacement surgery. MMP inhibitors are expected to halt or reverse the progression of cartilage loss and obviate or delay surgical intervention.\nProteases are critical elements at several stages in the progression of metastatic cancer. In this process, the proteolytic degradation of structural protein in the basal membrane allows for expansion of a tumor in the primary site, evasion from this site as well as homing and invasion in distant, secondary sites. Also, tumor induced angiogenesis is required for tumor growth and is dependent on proteolytic tissue remodeling. Transfection experiments with various types of proteases have shown that the matrix metalloproteases play a dominant role in these processes in particular gelatinases A and B (MMP-2 and MMP-9, respectively). For an overview of this field see Mullins, et al., Biochim. Biophys. Acta 695, 177 (1983); Ray, et al., Eur. Respir. J. 7, 2062 (1994); Birkedal-Hansen, et al., Crit. Rev. Oral Biol. Med. 4, 197 (1993).\nFurthermore, it was demonstrated that inhibition of degradation of extracellular matrix by the native matrix metalloprotease inhibitor TIMP-2 (a protein) arrests cancer growth (DeClerck, et al., Cancer Res. 52, 701 (1992)) and that TIMP-2 inhibits tumor-induced angiogenesis in experimental systems (Moses, et al. Science 248, 1408 (1990)). For a review, see DeClerck, et al., Ann. N. Y. Acad. Sci. 732, 222 (1994). It was further demonstrated that the synthetic matrix metalloprotease inhibitor batimastat when given intraperitoneally inhibits human colon tumor growth and spread in an orthotopic model in nude mice (Wang, et al. Cancer Res. 54, 4726 (1994)) and prolongs the survival of mice bearing human ovarian carcinoma xenografts (Davies, et. al., Cancer Res. 53, 2087 (1993)). The use of this and related compounds has been described in Brown, et al., WO-9321942 A2 (931111).\nThere are several patents and patent applications claiming the use of metalloproteinase inhibitors for the retardation of metastatic cancer, promoting tumor regression, inhibiting cancer cell proliferation, slowing or preventing cartilage loss associated with osteoarthritis or for treatment of other diseases as noted above (e.g. Levy, et al., WO-9519965 A1; Beckett, et al., WO-9519956 A1; Beckett, et al., WO-9519957 A1; Beckett, et al., WO-9519961 A1; Brown, et al., WO-9321942 A2; Crimmin, et al., WO-9421625 A1; Dickens, et al., U.S. Pat. No. 4,599,361; Hughes, et al., U.S. Pat. No. 5,190,937; Broadhurst, et al., EP 574758 A1; Broadhurst, et al,. EP 276436; and Myers, et al., EP 520573 A1. The preferred compounds of these patents have peptide backbones with a zinc complexing group (hydroxamic acid, thiol, carboxylic acid or phosphinic acid) at one end and a variety of sidechains, both those found in the natural amino acids as well as those with more novel functional groups. Such small peptides are often poorly absorbed, exhibiting low oral bioavailability. They are also subject to rapid proteolytic metabolism, thus having short half lives. As an example, batimastat, the compound described in Brown, et al., WO-9321942 A2, can only be given intra peritoneally.\nCertain 3-biphenoylpropanoic and 4-biaryloylbutanoic acids are described in the literature as anti-inflammatory, anti-platelet aggregation, anti-phlogistic, anti-proliferative, hypolipidemic, antirheumatic, analgesic, and hypocholesterolemic agents. In none of these examples is a reference made to MMP inhibition as a mechanism for the claimed therapeutic effect. Certain related compounds are also used as intermediates in the preparation of liquid crystals.\nSpecifically, Tomcufcik, et al., U.S. Pat. No. 3,784,701 claims certain substituted benzoylpropionic acids to treat inflammation and pain. These compounds include 3-biphenoylpropanoic acid (a.k.a fenbufen) shown below. \nChild, et al., J. Pharm. Sci., 66, 466 (1977) describes structure-activity relationships of several analogs of fenbufen. These include several compounds in which the biphenyl ring system is substituted or the propanoic acid portion is substituted with phenyl, halogen, hydroxyl or methyl, or the carboxylic acid or carbonyl functions are converted to a variety of derivatives. No compounds are described which contain a 4xe2x80x2-substituted biphenyl and a substituted propanoic acid portion combined in one molecule. The phenyl (compounds XLIX and LXXVII) and methyl (compound XLVII) substituted compounds shown below were described as inactive. \nKameo, et al., Chem. Pharm. Bull., 36, 2050 (1988) and Tomizawa, et al., JP patent 62132825 A2 describe certain substituted 3-biphenoylpropionic acid derivatives and analogs thereof including the following. Various compounds with other substituents on the propionic acid portion are described, but they do not contain biphenyl residues. \nwherein Xxe2x95x90H, 4xe2x80x2-Br, 4xe2x80x2-Cl, 4xe2x80x2-CH3, or 2xe2x80x2-Br.\nCousse, et al., Eur. J. Med. Chem., 22, 45 (1987) describe the following methyl and methylene substituted 3-biphenoyl-propanoic and -propenoic acids. The corresponding compounds in which the carbonyl is replaced with either CH2OH or CH2 are also described. \nwherein Xxe2x95x90H, Cl, Br, CH3O, F, or NH2.\nNichl, et al. DE patent 1957750 also describes certain of the above methylene substituted biphenoylpropanoic acids.\nEl-Hashash, et al., Revue Roum. Chim. 23, 1581 (1978) describe products derived from xcex2-aroyl-acrylic acid epoxides including the following biphenyl compound. No compounds substituted on the biphenyl portion are described. \nKitamura, et al., JP patent 60209539 describes certain biphenyl compounds used as intermediates for the production of liquid crystals including the following. The biphenyl is not substituted in these intermediates. \nwherein R1 is an allyl of 1-10 carbons.\nThyes, et al., DE patent 2854475 uses the following compound as an intermediate. The biphenyl group is not substituted. \nSammour, et al., Egypt J. Chem. 15, 311 (1972) and Couquelet, et al., Bull. Soc. Chim. Fr. 9, 3196 (1971) describe certain dialkylamino substituted biphenoylpropanoic acids including the following. In no case is the biphenyl group substituted. \nwherein R1, R2=alkyl, benzyl, H, or, together with the nitrogen, morpholinyl.\nOthers have disclosed a series of biphenyl-containing carboxylic acids, illustrated by the compound shown below, which inhibit neural endopeptidase (NEP 24.11), a membrane-bound zinc metalloprotease (Stanton, et al., Bioorg. Med. Chem. Lett. 4, 539 (1994); Lombaert, et al., Bioorg. Med. Chem. Lett. 4, 2715 (1994); Lombaert, et al., Bioorg. Med. Chem. Lett. 5, 145 (1995); Lombaert, et al., Bioorg. Med. Chem. Lett. 5, 151 (1995)). \nIt has been reported that N-carboxyalkyl derivatives containing a biphenylethylglycine, illustrated by the compound shown below, are inhibitors of stromelysin-1 (MMP-3), 72 kDA gelatinase (MMP-2) and collagenase (Durette, et al., WO-9529689). \nIt would be desirable to have effective MMP inhibitors which possess improved bioavailability and biological stability relative to the peptide-based compounds of the prior art, and which can be optimized for use against particular target MMPs. Such compounds are the subject of the present application.\nThe development of efficacious MMP inhibitors would afford new therapies for diseases mediated by the presence of, or an excess of MMP activity, including osteoarthritis, rheumatoid arthritis, septic arthritis, tumor metastasis, periodontal diseases, corneal ulcerations, and proteinuria. Several inhibitors of MMPs have been described in the literature, including thiols (Beszant, et al., J. Med. Chem. 36, 4030 (1993), hydroxamic acids (Wahl, et al. Bioorg. Med. Chem. Lett. 5, 349 (1995) Conway, et al. J. Exp. Med. 182, 449 (1995); Porter, et al., Bioorg. Med. Chem. Lett. 4, 2741 (1994); Tomczuk, et al., Bioorg. Med. Chem. Lett. 5, 343 (1995); Castelhano, et al., Bioorg. Med. Chem. Lett. 5, 1415 (1995)), phosphorous-based acids (Bird, et al. J. Med. Chem. 37, 158 (1994); Morphy, et al., Bioorg. Med. Chem. Lett. 4, 2747 (1994); Kortylewicz, et al., J. Med. Chem. 33, 263 (1990)), and carboxylic acids (Chapman, et al. J. Med. Chem. 36, 4293 (1993); Brown, et al. J. Med. Chem. 37, 674 (1994); Morphy, et al., Bioorg. Med. Chem. Lett. 4, 2747 (1994); Stack, et al., Arch. Biochem. Biophys. 287, 240 (1991); Ye, et al., J. Med. Chem. 37, 206 (1994); Grobelny, et al., Biochemistry 24, 6145 (1985); Mookhtiar, et al., Biochemistry 27, 4299 (1988)). However, these inhibitors generally contain peptidic backbones, and thus usually exhibit low oral bioactivity due to poor absorption and short half lives due to rapid proteolysis. Therefore, there remains a need for improved MMP inhibitors.\nThis invention provides compounds having matrix metalloprotease inhibitory activity. These compounds are useful for inhibiting matrix metalloproteases and, therefore, combating conditions to which MMP\"\"s contribute. Accordingly, the present invention also provides pharmaceutical compositions and methods for treating such conditions.\nThe compounds described relate to a method of treating a mammal comprising administering to the mammal a matrix metalloprotease inhibiting amount of a compound according to the invention sufficient to:\n(a) alleviate the effects of osteoarthritis, rheumatoid arthritis, septic arthritis, periodontal disease, corneal ulceration, proteinuria, aneurysmal aortic disease, dystrophobic epidermolysis, bullosa, conditions leading to inflammatory responses, osteopenias mediated by MMP activity, tempero mandibular joint disease, demyelating diseases of the nervous system;\n(b) retard tumor metastasis or degenerative cartilage loss following traumatic joint injury;\n(c) reduce coronary thrombosis from athrosclerotic plaque rupture; or\n(d) temporarily reduce fertility (i.e., act as effective birth control agents).\nThe compounds of the present invention are also useful scientific research tools for studying functions and mechanisms of action of matrix metalloproteases in both in vivo and in vitro systems. Because of their MMP-inhibiting activity, the present compounds can be used to modulate MSP action, thereby allowing the researcher to observe the effects of reduced MMP activity in the experimental biological system under study.\nThis invention relates to compounds having matrix metalloprotease inhibitory activity and the generalized formula:\n(T)xA-B-D-E-Gxe2x80x83xe2x80x83(L)\nIn the above generalized formula (L), (T)xA represents a substituted or unsubstituted aromatic 6-membered ring or heteroaromatic 5-6 membered ring containing 1-2 atoms of N, O, or S. T represents one or more substituent groups, the subscript x represents the number of such substuent groups, and A represents the aromatic or heteroaromatic ring.\nIn the above generalized formula (L), TxA represents a substituted or unsubstituted aromatic 6-membered ring or heteroaromatic 5-6 membered ring containing 1-2 atoms independently selected from the group of N, O, or S. T represents a substituted acetylenic moiety.\nIn the generalized formula (L), B represents an aromatic 6-membered ring or a heteroaromatic 5-6 membered ring containing 1-2 atoms independently selected from the group of N, O, or S. It is referred to as the B ring or B unit. When N is employed in conjunction with either S or O in the B ring, these heteroatoms are separated by at least one carbon atom.\nIn the generalized formula (L), D represents \nIn the generalized formula (L), E represents a chain of n carbon atoms bearing m substituents R6 in which the R6 groups are independent substituents, or constitute spiro or nonspiro rings. Rings may be formed in two ways: a) two groups R6 are joined, and taken together with the chain atom(s) to which the two R6 group(s) are attached, and any intervening chain atoms, constitute a 3-7 membered ring, or b) one group R6 is joined to the chain on which this one group R6 resides, and taken together with the chain atom(s) to which the R6 group is attached, and any intervening chain atoms, constitutes a 3-7 membered ring. The number n of carbon atoms in the chain is 2 or 3, and the number m of R6 substituents is an integer of 1-3. The number of carbons in the totality of R6 groups is at least two.\nEach group R6 is alkyl, alkenyl, alkynyl, heteroaryl, non-aromatic cyclic, and combinations thereof optionally substituted with one or more heteroatoms\nIn the generalized formula (L), E preferably represents a linear or cyclic alkyl moiety substituted with a mono- or -biheterocyclic ring structure.\nIn the generalized formula (L), G represents xe2x80x94PO3H2, xe2x80x94M, \nin which M represents xe2x80x94CO2H, xe2x80x94CON(R11)2 wherein R11 is H or simple alkyl, or xe2x80x94CO2R12 wherein R12 is lower alkyl, and R13 represents any of the side chains of the 19 noncyclic naturally occurring amino acids.\nPharmaceutically acceptable salts of these compounds are also within the scope of the invention.\nIn most related reference compounds of the prior art, the biphenyl portion of the molecule is unsubstituted, and the propanoic or butanoic acid portion is either unsubstituted or has a single methyl or phenyl group. Presence of the larger phenyl group has been reported to cause prior art compounds to be inactive as anti-inflammatory analgesic agents. See, for example, Child, et al., J. Pharm. Sci. 66, 466(1977). By contrast, it has now been found that compounds which exhibit potent MMP inhibitory activity contain a substituent of significant size on the propanoic or butanoic portion of the molecule. The biphenyl portions of the best MMP inhibitors also preferably contain a substituent on the 4xe2x80x2-position, although when the propanoic or butanoic portions are optimally substituted, the unsubstituted biphenyl compounds of the invention have sufficient activity to be considered realistic drug candidates.\nThe foregoing merely summarizes certain aspects of the present invention and is not intended, nor should it be construed, to limit the invention in any way. All of the patents and other publications recited in this specification are hereby incorporated by reference in their entirety.\nMore particularly, the compounds of the present invention are materials having matrix metalloprotease inhibitory activity and the generalized formula:\n(T)xA-B-D-E-Gxe2x80x83xe2x80x83(L)\nin which (T)xA represents a substituted or unsubstituted aromatic or heteroaromatic moiety selected from the group consisting of: \nin which R1 represents H or alkyl of 1-3 carbons.\nThroughout this application, in the displayed chemical structures, an open bond indicates the point at which the structure joins to another group. For example, \nIn the above structures for (T)xA, the aromatic ring is referred to as the A ring or A unit, and T represents a substituent group, referred to as a T group or T unit. x is preferably 1.\nThe B ring of generalized formula (L) is a substituted or unsubstituted aromatic or heteroaromatic ring, in which any substituents are groups which do not cause the molecule to fail to fit the active site of the target enzyme, or disrupt the relative conformations of the A and B rings, such that they would be detrimental. Such substituents may be moieties such as lower alkyl, lower alkoxy, CN, NO2, halogen, etc., but are not to be limited to such groups.\nIn the generalized formula (L), B represents an aromatic or heteroaromatic ring selected from the group consisting of: portion comprises 4-9 carbons and at least one N, O, or S heteroatom and the alkyl portion contains 1-4 carbons.\nR4 represents H; alkyl of 1-12 carbons; aryl of 6-10 carbons; heteroaryl comprising 4-9 carbons and at least one N, O, or S heteroatom; arylalkyl in which the aryl portion contains 6-10 carbons and the alkyl portion contains 1-4 carbons; heteroaryl-alkyl in which the heteroaryl portion comprises 4-9 carbons and at least one N, O, or S heteroatom and the alkyl portion contains 1-4 carbons; alkenyl of 2-12 carbons; alkynyl of 2-12 carbons; xe2x80x94(CqH2qO)rR5 in which q is 1-3, r is 1-3, and R5 is H provided q is greater than 1, or R5 is alkyl of 1-4 carbons, or phenyl; xe2x80x94(CH2)sX in which s is 2-3 and X is halogen; or xe2x80x94C(O)R2.\nAny unsaturation in a moiety which is attached to Q or which is part of Q is separated from any N, O, or S of Q by at least one carbon atom, and the number of substituents, designated x, is 0, 1, or 2.\nThe substituent group T can be an acetylene containing moiety with the general formula:\nR30(CH2)nxe2x80x2Cxe2x89xa1Cxe2x80x94\nwhere n is 1-4 and R30 is selected from the group consisting of: HOxe2x80x94, MeOxe2x80x94, N(n-Pr)2xe2x80x94, CH3CO2xe2x80x94, CH3CH2OCO2xe2x80x94, HO2Cxe2x80x94, OHCxe2x80x94, Phxe2x80x94, 3-HOxe2x80x94Phxe2x80x94, and PhCH2Oxe2x80x94, provided that when R30 is Ph or 3-HOxe2x80x94Ph, n=0.\nThe B ring of generalized formula (L) is a substituted or unsubstituted aromatic or heteroaromatic ring, in which any substituents are groups which do not cause the molecule to fail to fit the active site of the target enzyme, or disrupt the relative conformations of the A and B rings, such that they would be detrimental. Such substituents may be moieties such as lower alkyl, lower alkoxy, CN, NO2 , halogen, etc., but are not to be limited to such groups.\nIn the generalized formula (L), B represents an aromatic or heteroaromatic ring selected from the group consisting of: \nin which R1 is defined as above. These rings are referred to as the B ring or B unit.\nIn an alternative embodiment, compounds of the general formula (L) include those in which the combination (T)xxe2x80x94Axe2x80x94B has the structure: \nwhere Z may be (CH2)exe2x80x94C6H4xe2x80x94(CH2)f or (CH2)g, e=0-8, f=0-5 and g=0-14, rxe2x80x2 is 0-6. R15 may be a straight, or cyclic alkyl group of 6-12 carbons atoms, preferably of 7-11 carbon atoms, and optionally may bear one or more pharmaceutically acceptable substituents which are discussed more fully below.\nR15 may also be a polyether of the formula R32O(C2H4O)h in which the subscript xe2x80x9chxe2x80x9d is 1 or 2, and the group R32 is a straight, branched or cyclic alkyl group of 1-5 carbon atoms, preferably of 1-3 carbon atoms and straight, or phenyl, or benzyl. R32 optionally may bear one or more pharmaceutically-acceptable substituents.\nR15 may also be a substituted alkynyl group of the formula:\nR33(CH2)bxe2x80x94Cxe2x89xa1Cxe2x80x94\nin which the subscript xe2x80x9cbxe2x80x9d is 1-10 and the group R33 is Hxe2x80x94, HOxe2x80x94 or R34Oxe2x80x94 and the group is preferably the HOxe2x80x94 group. R34 may be an alkyl group of 1-3 carbon atoms, or phenyl or benzyl. R33 optionally may bear one or more pharmaceutically-acceptable substituents.\nR15 may also be H, Cl, MeO or \nwherein n is 0-4, R17 is C2H5, allyl, or benzyl.\nIn the generalized formula (L), D represents the moieties: \nIn the generalized formula (L), E represents a moiety having between D and G shown by the following formula: \nwherein r is 0-2 and R40 is a mono- or bi- heterocyclic structure. When r=0 the above structure takes the form \nWhen r is 1 or 2, a cyclobutyl or cyclopentyl ring is formed, respectively. Each ring of the mono- or bi- heterocylic structures comprise 5-7 membered rings substituted with 1-3 heteroatoms independently selected from N, S, and O; one or two carbons of the ring are optionally carbonyl carbons; any sulfur of the ring is optionally xe2x80x94S(O)xe2x80x94 or xe2x80x94S(O)2xe2x80x94; one or more ring members are optionally substituted with one or two methyl groups.\nIn addition, aryl or heteroaryl portions of any of the T or R6 groups optionally may bear up to two substituents such as xe2x80x94(CH2)yC(R11)(R12)OH, xe2x80x94(CH2)yOR11, xe2x80x94(CH2)ySR11, xe2x80x94(CH2)yS(O)R11, xe2x80x94(CH2)yS(O)2R11, xe2x80x94(CH2)ySO2N(R11)2, xe2x80x94(CH2)yN(R11)2, xe2x80x94(CH2)yN(R11)COR12, xe2x80x94OC(R11)2Oxe2x80x94 in which both oxygen atoms are connected to the aryl ring.\nThe B ring is preferably a 1,4-phenylene or 2,5-thiophene ring, most preferably 1,4-phenylene.\nThe D unit is most preferably a carbonyl group.\nIn the E unit, r is preferably 0 or 2 and R40 is preferably one of the following: \nor PhCH2OCH2OCH2xe2x80x94.\nThe G unit is most preferably a carboxylic acid group and is attached to the E unit at the 2 position, i.e., the carbon atom of the E unit beta to the D unit.\nIt is to be understood that as used herein, the term xe2x80x9calkylxe2x80x9d means straight, branched cyclic, and polycyclic materials. The term xe2x80x9chaloalkylxe2x80x9d means partially or fully halogenated alkyl groups such as xe2x80x94(CH2)2Cl, xe2x80x94CF3 and xe2x80x94C6F13, for example.\nIn the generalized formula (L), the A and B rings are preferably phenyl and phenylene, respectively, the A ring preferably bears at least one substiuent group T preferably located on the position furthest from the position of the A ring which is connected to the B ring, the D unit is preferably a carbonyl group, and the G unit is preferably a carboxyl group.\nCertain alternative embodiments include compounds having matrix metalloproteinase inhibitory activity and the following generalized formula: \nwhere Z=(CH2)exe2x80x94C6H4xe2x80x94(CH2)f or (CH2)g, e=0-8, f=0-5, g=0-14, rxe2x80x2 is 0-6 and where y is 0, 2, or 3.\nR15 may be H, Cl, MeO or \nwherein nxe2x80x3 is 0-4, R17 is C2H5, allyl or benzyl, and R40 is one of: \nand xe2x80x94CH2OCH2OCH2Ph.\nThe most preferred compounds of generalized formula (L) are \nwherein T is selected from a group consisting of: \nr is 0-2, and R40 is selected from the group consisting of: \nand xe2x80x94CH2OCH2OCH2Ph.\nThe invention also relates to certain intermediates useful in the synthesis of some of the claimed inhibitors. These intermediates are compounds having the generalized formula \nwhere Bn is benzyl, TMSE is trimethylsilyl ethyl and R40 is as defined above.\nThose skilled in the art will appreciate that many of the compounds of the invention exist in enantiomeric or diastereomeric forms, and that it is understood by the art that such stereoisomers generally exhibit different activities in biological systems. This invention encompasses all possible stereoisomers which possess inhibitory activity against an MMP, regardless of their stereoisomeric designations, as well as mixtures of stereoisomers in which at least one member possesses inhibitory activity.\nThe most preferred compounds of the present invention are as indicated and named in the list below:\nI) 2-[(4xe2x80x2-chloro [1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(4-oxo-1,2,3-benzotriazin-3(4H)-yl)methyl]-cyclopentanecarboxylic acid,\nII) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[phenoxymethoxymethyl]-cyclopentanecarboxylic acid,\nIII) 2-[4xe2x80x2-chloro [1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[[(1-pyrrolidinylthioxomethyl)thio]methyl]-cyclopentanecarboxylic acid,\nIV) 2-[(4xe2x80x2-choro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(1,1-dioxido-3-oxo-1,2-benzisothiazol-2(3H)-yl)methyl]-cyclopentanecarboxylic acid,\nV) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[1-oxo-2(1H)-phthalazinyl)methyl]-cyclopentanecarboxylic acid,\nVI) 2-[4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(2-oxo-3(2H)-benzoxazolyl)methyl]-cyclopentanecarboxylic acid,\nVII) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[5,5-dimethyl-2,4-dioxo-3-oxazolidinyl-methyl]-cyclopentanecarboxylic acid,\nVIII) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(2,4-dioxo-3-thiazolidinyl)methyl]-cyclopentanecarboxylic acid,\nIX) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]4-yl)carbonyl]-5-[2,4,5-trioxo-1-imidazolidinyl)methyl]-cyclopentanecarboxyl acid\nX) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(3,6-dihydro-2,6-dioxo-1(2H)-pyrimidinyl)methyl]-cyclopentanecarboxylic acid,\nXI) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[3,4-dihydro-2,4-dioxo-1(2H)-pyrimidinyl)methyl]-cyclopentanecarboxylic acid,\nXII) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(1,4-dihydro-2,4-dioxo-3(2H)-quinazolinyl)methyl]-cyclopentanecarboxylic acid,\nXIII) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[3,4-dihydro-1,3-dioxo-2(1H)-isoquinolinyl)methyl]-cyclopentanecarboxylic acid,\nXIV) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(1,4-dihydro-4-oxo-3(2H)-quinazolinyl)methyl]-clopentanecarboxylic acid,\nXV) 2-[4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]4-yl)carbonyl]-5-[(1,3-dihydro-3-oxo-2H-indazol-2-yl)methyl]-cyclopentanecarboxylic acid,\nXVI) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[2,3-dihydro-1H-benzimidazol-1-yl)methyl]-cyclopentanecarboxylic acid,\nXVII) 2-[(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)carbonyl]-5-[(3,4-dihydro-1,4-dioxo-2(1H)-phthalazinyl)methyl]-clopentanecarboxylic acid,\nXVIII) R/S xcex1-[2-(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)-2-oxoethyl]-1-oxo-2(1H)-phthalazinebutanoic acid,\nXIX) R-xcex1-[2-(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]4-yl)-2-oxoethyl]-1-oxo-2(1H)-phthalazinebutanoic acid,\nXX) S-xcex1-[2-(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)-2-oxoethyl]-1-oxo-2(1H)-phthalazinebutanoic acid,\nXXI) xcex1-[2-(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)-2-oxoethyl]-4-oxo-1,2,3,-Benzotriazine-3(4H)-butanoic acid, and\nXXII) xcex1-[2-(4xe2x80x2-chloro[1,1xe2x80x2-biphenyl]-4-yl)-2-oxoethyl]-2,3-dihydro-5-methyl-2-oxo-1H-1,4-benzodiazepine-1-butanoic acid.\nGeneral Preparative Methods:\nThe compounds of the invention may be prepared readily by use of known chemical reactions and procedures. Nevertheless, the following general preparative methods are presented to aid the reader in synthesizing the inhibitors, with more detailed particular examples being presented below in the experimental section describing the working examples. All variable groups of these methods are as described in the generic description if they are not specifically defined below. The variable subscript n is independently defined for each method. When a variable group with a given symbol (i.e R9) is used more than once in a given structure, it is to be understood that each of these groups may be independently varied within the range of definitions for that symbol.\nGeneral Method Axe2x80x94The compounds of this invention in which the rings A and B are substituted phenyl and phenylene respectively are conveniently prepared by use of a Friedel-Crafts reaction of a substituted biphenyl MII with an activated acyl-containing intermediate such as the succinic or glutaric anhydride derivative MIII or acid chloride MIV in the presence of a Lewis acid catalyst such as aluminum trichloride in an aprotic solvent such as 1,1,2,2-tetrachloroethane. The well known Friedel-Crafts reaction can be accomplished with use of many alternative solvents and acid catalysts as described by Berliner, Org. React, 5, 229, 1949 and Heaney, Comp. Org. Synth. 2, 733, 1991.\nIf the anhydride MIII is monosubstituted or multiply-substituted in an unsymmetrical way, the raw product MI-A often exists as a mixture of isomers via attack of the anhydride from either of the two carbonyls. The resultant isomers can be separated into pure forms by crystallization or chromatography using standard methods known to those skilled in the art.\nWhen they are not commercially available, the succinic anhydrides MIII can be prepared via a Stobbe Condensation of a dialkyl succinate with an aldehyde or ketone (resulting in side chain R6), followed by catalytic hydrogenation, hydrolysis of a hemiester intermediate to a diacid, and then conversion to the anhydride MIII by reaction with acetyl chloride or acetic anhydride. Alternatively, the hemiester intermediate is converted by treatment with thionyl chloride or oxalyl chloride to the acid chloride MIV. For a review of the Stobbe condensation, including lists of suitable solvents and bases see Johnson and Daub, Org. React., 6, 1, 1951.\nThis method, as applied to the preparation of MIII (R6xe2x95x90H, isobutyl and H, n-pentyl), has been described Wolanin, et al., U.S. Pat. No. 4,771,038. \nMethod A is especially useful for the preparation of cyclic compounds such as MI-A-3, in which two R6 groups are connected in a methylene chain to form a 3-7 member ring. Small ring (3-5 member) anhydrides are readily available only as cis isomers which yield cis invention compounds MI-A-3. The trans compounds MI-A-4 are then prepared by treatment of MI-A-3 with a base such as DBU in THF. The substituted four member ring starting material anhydrides such as MIII-A-1 are formed in a photochemical 2+2 reaction as shown below. This method is especially useful for the preparation of compounds in which R40 is acetoxy or acetoxymethylene. After the subsequent Friedel-Crafts reaction the acetate can be removed by basic hydrolysis and the carboxyl protected by conversion to 2-(trimethylsilyl)ethyl ester. The resultant intermediate with R40=CH2OH can be converted to invention compounds with other R40 groups by using procedures described in General Method G. \nThe Friedel-Crafts method is also useful when double bonds are found either between C-2 and C-3 of a succinoyl chain (from maleic anhydride or 1-cyclopentene-1,2-dicarboxylic anhydride, for example) or when a double bond is found in a side chain, such as in the use of itaconic anhydride as stating material to yield products in which two R6 groups are found on one chain carbon together to form an exo-methylene (xe2x95x90CH2) group. Subsequent uses of these compounds are described in Methods D.\nGeneral Method Bxe2x80x94Alternatively the compounds MI can be prepared via a reaction sequence involving mono-alkylation of a dialkyl malonate MVI with an alkyl halide to form intermediate MVII, followed by alkylation with a halomethyl biphenyl ketone MVIII to yield intermediate MIX. Compounds of structure MIX are then hydrolyzed with aqueous base and heated to decarboxylate the malonic acid intermediate and yield MI-B-2 (Method B-1). By using one equivalent of aqueous base the esters MI-B-2 with R12 as alkyl are obtained, and using more than two equivalents of base the acid compounds (R12=H) are obtained. Optionally, heat is not used and the diacid or acid-ester MI-B-1 is obtained.\nAlternatively, the diester intermediate MIX can be heated with a strong acids such as concentrated hydrochloric acid in acetic acid in a sealed tube at about 110xc2x0 C. for about 24 hr to yield MI-B-1 (R12=H). Alternatively, the reaction of MVI with MVIII can be conducted before that with the alkyl halide to yield the same MIX (Method B-2).\nAlternatively, a diester intermediate MXIX, which contains R12=allyl, can be exposed to Pd catalysts in the presence of pyrrolidine to yield MI-B-2 (R12=H) (Dezeil, Tetrahedron Lett. 28, 4371, 1990.\nIntermediates MVII are formed from biphenyis MII in a Friedel-Craft reaction with haloacetyl halides such as bromoacetyl bromide or chloroacetyl chloride. Alternatively, the biphenyl can be reacted with acetyl chloride or acetic anhydride and the resultant product halogenated with, for example, bromine to yield intermediates MVIII (X=Br).\nMethod B has the advantage of yielding single region isomers when Method A yields mixtures. Method B is especially useful when the side chains R6 contain aromatic or heteroaromatic rings that may participate in intramolecular acylation reactions to give side products if Method A were to be used. This method is also very useful when the R6 group adjacent to the carboxyl of the final compound contains heteroatoms such as oxygen, sulfur, or nitrogen, or more complex functions such as imide rings. \nWhen R6 contains selected functional groups Z, malonate MVII can be prepared by alkylating a commercially available unsubstituted malonate with prenyl or allyl halide, subject this product to ozonalysis with reductive work-up, and the desired z group can be coupled via a Mitsunobu reaction (Mitsunobu, Synthesis 1, 1981). Alternatively, the intermediate alcohol can be subjected to alkylation conditions to provide malonate MVII containing the desired Z group. \nGeneral Method Cxe2x80x94Especially useful is the use of chiral HPLC to separate the enantiomers of racemic product mixtures (see, for example, Arit, et al., Chem. Int. Ed. Engl. 12, 30 (1991)). The compounds of this invention can be prepared as pure enantiomers by use of a chial auxiliary route. See, for example, Evans, Aldrichimica Acta, 15(2), 23, 1982 and other similar references known to one skilled in the art.\nGeneral Method Dxe2x80x94Compounds in which R6 are alkyl- or aryl- or heteroaryl- or acyl- or heteroarylcarbonyl-thiomethylene are prepared by methods analogous to those described in the patent WO 90/05719. Thus substituted itaconic anhydride MXVI (n=1) is reacted under Friedel-Crafts conditions to yield acid MI-D-1 which can be separated by chromatography or crystallization from small amounts of isomeric MI-D-5. Alteratively, MI-D-5s are obtained by reaction of invention compounds MI-D4 (from any of Methods A through C) with formaldehyde in the presence of base.\nCompounds MI-D-1 or MI-D-5 are then reacted with a mercapto derivative MXVII or MXVIII in the presence of catalyst such as potassium carbonate, ethyldiisobutylamine, tetrabutylammonium fluoride or free radical initiators such as azobisisobutyronitrile (AIBN) in a solvent such as diethylformamide or tetrahydrofuiran to yield invention compounds MI-D-2, MI-D-3, MI-D-6, or MI-D-7. \nGeneral Method Exe2x80x94Biaryl compounds such as those of this application may also be prepared by Suzuki or Stille cross-coupling reactions of aryl or heteroaryl metallic compounds in which the metal is zinc, tin, magnesium, lithium, boron, silicon, copper, cadmium or the like with an and or heteroaryl halide or triflate (trifluoromethane-sulfonate) or the like. In the equation below either Met or X is the metal and the other is the halide or triflate (OTf). Pd(com) is a soluble complex of palladium such as tetrakis(triphenylphosphine)-palladium(O) or bis- (triphenylphosphine)-palladium(III) chloride. These methods are well known to those skilied in the art. See, for example, Suzuki, Pure Appl. Chem.63, 213 (1994); Suzuki, Pure Appl. Chem. 63, 419 (1991); and Farina and Roth, xe2x80x9cMetal-Organic Chemistryxe2x80x9d Volume 5 (Chapter 1), 1994.\nThe starting materials MXXIII (B=1,4-phenylene) are readily formed using methods analogous to those of methods A, B, C, or D but using a halobenzene rather than a biphenyl as starting material. When desired, the materials in which X is halo can be converted to those in which X is metal by reactions well known to those skilled in the art, such as treatment of a bromo intermediate with hexamethylditin and palladium tetrakistriphenylphosphine in toluene at reflux to yield the trimethyltin intermediate. The starting materials MXXIII (B=heteroaryl) are most conveniently prepared by method C but using readily available heteroaryl rather than biphenyl starting materials. The intermediates MXXII are either commercial or easily prepared from commercial materials by methods well known to those skilled in the art. \nT, x, A, B, E and G as in Structure (L)\nMet=Metal and X=Halide or Triflate or\nMet=Halide or Triflate arid X=Metal\nThese general methods are useful for the preparation of compounds for which Friedel-Crafts reactions such as those of Methods A, B, C, or D would lead to mixtures with various biaryl acylation patterns. Method E is also especially useful for the preparation of products in which the aryl groups, A or B, contain one or more heteroatoms (heteroaryls) such as those compounds that contain thiophene, fuiran, pyridine, pyrrole, oxazole, thiazole, pyriridine or pyrazine rings or the like instead of phenyls.\nGeneral Method Fxe2x80x94When the R6 groups of method F form together a 4-7 member carbocyclic ring as in Intermediate MXXV below, the double bond can be moved out of conjugation with the ketone group by treatment with two equivalents of a strong base such as lithium diisopropylamide or lithium hexamethylsilylamide or the like followed by acid quench to yield compounds with the structure MXXVI. Reaction of MXXVI with mercapto derivatives using methods analogous to those of General Method D then leads to cyclic compounds MI-F-I or MI-F-2. \nGeneral Method Gxe2x80x94The compounds of this invention in which two R6 groups are joined to form a substituted 5-member ring are most conveniently prepared by method G. In this method acid CLII (R=H) is prepared using the protocols described in Tetrahedron 37, Suppl., 411 (1981). The acid is protected as an ester [eg. R=benzyl (Bn) or 2-(trimethylsilyl)ethyl TMSE)] by use of coupling agents such as 1-(3-dimethylaminopropyl)-3-ethylcarbodiimide hydrochloride and procedures well known to those skilled in the art. Substituted bromobiphenyl CIII is converted to its Grignard reagent by treatment with magnesium and reacted with CLII to yield alcohol CVI. Alcohol CVI is eliminated via base treatment of its mesylate by using conditions well known to those skilled in the art to yield olefin CVII. Alternatively CIII is converted to a trimethyltin intermediate via initial metallation of the bromide with n-butyllithium at low temperature (xe2x88x9278xc2x0 C.) followed by treatment with chlorotrimethyltin and CI is converted to an enoltriflate (CII) by reaction with 2-[N,N-bis(trifluoromethylsulfonyl)amino]-5-chloropyridine in the presence of a strong aprotic base. The tin and enoltriflate intermediates are then coupled in the presence of a Pd0 catalyst, CuI and AsPh3 to yield directly intermediate CVII. Ozonolysis of CVII (workup with methylsufide) yields aldehyde CVIII. Alternatively treatment with OsO4 followed by HIO4 converts CVII to CVIII. \nConversion of key intermediate CVIII to the targeted patent compounds is accomplished in several ways depending on the identity of side chain function Z. Reaction of CVIII with Wittig reagents followed by hydrogenation yields products in which Z is alkyl and or arylalkyl. Selective reduction of aldehyde CVIII with a reducing agent such as lithium tris [(3-ethyl-3pentyl)oxy]aluminum hydride (LTEPA) yields alcohol CIX. The alcohol is converted to phenyl ethers or a variety of heteroatom substituted derivatives used to generate sidechain Z via the Mitsunobu reaction using conditions well known to those skilled in the art (see Mitsunobu, Synthesis, 1 (1981)). Alternatively the alcohol of CIX is converted to a leaving group such as tosylate (CX) or bromide by conditions well known to those skilled in the art and then the leaving group is displaced by an appropriate nucleophile. Several examples of this type of reaction can be found in Norman, et al., J. Med. Chem. 37, 2552 (1994). Direct acylation of the alcohol CIX yields invention compounds in which Z=OAcyl and reaction of the alcohol with various alkyl halides in the presence of base yields alkyl ethers. In each case a final step is removal of acid blocking group R to yield acids (R=H) by using conditions which depend on the stability of R and Z, but in all cases well known to those skilled in the art such as removal of benzyl by base hydrolysis or of 2-(trimethylsilyl)ethyl by treatment with tetrabutylammonium fluoride.\nGeneral Method Hxe2x80x94Amides of the acids of the invention compounds can be prepared from the acids by treatment in an appropriate solvent such as dichloromethane or dimethylformamide with a primary or secondary amine and a coupling agent such as dicyclohexylcarbodiimide. These reactions are well known to those skilled in the art. The amine component can be simple alkyl or arylalkyl substituted or can be amino acid derivatives in which the carboxyl is blocked and the amino group is free.\nGeneral Method Ixe2x80x94The compounds of this invention in which (T)x is an alkynyl or substituted alkynyl are prepared according to general method I (Austin, J. Org. Chem. 46, 2280 (1981)). Intermediate MX is prepared according to methods A, B, C, D or G by starting with commercial MIII (T=Br). Reaction of MX with substituted acetylene MXI in the presence of Cu(I)/palladate reagent gives invention compound MI-I-1. In certain cases, R3 may be an alcohol blocked as trialkylsilyl. In such cases the silyl group can be removed by treatment with acids such as trifluoroacetic acid or HFxe2x80x94pyridine reagent. \nSuitable pharmaceutically acceptable salts of the compounds of the present invention include addition salts formed with organic or inorganic bases. The salt forming ion derived from such bases can be metal ions, e.g., aluminum, alkali metal ions, such as sodium of potassium, alkaline earth metal ions such as calcium or magnesium, or an amine salt ion, of which a number are known for this purpose. Examples include ammonium salts, arylalkylamines such as dibenzylamine and N,N-dibenzylethylenediamine, lower alkylamines such as methylamine, t-butylamine, procaine, lower alkylpiperidines such as N-ethylpiperidine, cycloalkylamines such as cyclohexylamine or dicyclohexylamine, 1-adamantylamine, benzathine, or salts derived from amino acids like arginine, lysine or the like. The physiologically acceptable salts such as the sodium or potassium salts and the amino acid salts can be used medicinally as described below and are preferred.\nThese and other salts which are not necessarily physiologically acceptable are useful in isolating or purifying a product acceptable for the purposes described below. For example, the use of commercially available enantiomerically pure amines such as (+)-cinchonine in suitable solvents can yield salt crystals of a single enatiomer of the invention compounds, leaving the opposite enantiomer in solution in a process often referred to as xe2x80x9cclassical resolution.xe2x80x9d As one enantiomer of a given invention compound is usually substantially greater in physiological effect than its antipode, this active isomer can thus be found purified in either the crystals or the liquid phase. The salts are produced by reacting the acid form of the invention compound with an equivalent of the base supplying the desired basic ion in a medium in which the salt precipitates or in aqueous medium and then lyophilizing. The free acid form can be obtained from the salt by conventional neutralization techniques, e.g., with potassium bisulfate, hydrochloric acid, etc.\nThe compounds of the present invention have been found to inhibit the matrix metalloproteases MMP-3, MMP-9 and MMP-2. and to a lesser extent MMP-1, and are therefore useful for treating or preventing the conditions referred to in the background section. As other MMPs not listed above share a high degree of homology with those listed above, especially in the catalytic site, it is deemed that compounds of the invention should also inhibit such other MMPs to varying degrees. Varying the substituents on the biaryl portions of the molecules, as well as those of the propanoic or butanoic acid chains of the claimed compounds, has been demonstrated to affect the relative inhibition of the listed MMPs. Thus compounds of this general class can be xe2x80x9ctunedxe2x80x9d by selecting specific substituents such that inhibition of specific MMP(s) associated with specific pathological conditions can be enhanced while leaving non-involved MMPs less affected.\nThe method of treating matrix metalloprotease-mediated conditions may be practiced in mammals, including humans, which exhibit such conditions.\nThe inhibitors of the present invention are contemplated for use in veterinary and human applications. For such purposes, they will be employed in pharmaceutical compositions containing active ingredient(s) plus one or more pharmaceutically acceptable carriers, diluents, fillers, binders, and other excipients, depending on the administration mode and dosage form contemplated.\nAdministration of the inhibitors may be by any suitable mode known to those skilled in the art. Examples of suitable parenteral administration include intravenous, intraarticular, subcutaneous and intramuscular routes. Intravenous administration can be used to obtain acute regulation of peak plasma concentrations of the drug. Improved half-life and targeting of the drug to the joint cavities may be aided by entrapment of the drug in liposomes. It may be possible to improve the selectivity of liposomal targeting to the joint cavities by incorporation of ligands into the outside of the liposomes that bind to synovial-specific macromolecules. Alternatively intramuscular, intraarticular or subcutaneous depot injection with or without encapsulation of the drug into degradable microspheres e.g., comprising poly(DL-lactide-co-glycolide) may be used to obtain prolonged sustained drug release. For improved convenience of the dosage form it may be possible to use an i.p. implanted reservoir and septum such as the Percuseal system available from Pharmacia. Improved convenience and patient compliance may also be achieved by the use of either injector pens (e.g. the Novo Pin or Q-pen) or needle-free jet injectors (e.g. from Bioject, Mediject or Becton Dickinson). Prolonged zero-order or other precisely controlled release such as pulsatile release can also be achieved as needed using implantable pumps with delivery of the drug through a cannula into the synovial spaces. Examples include the subcutaneously implanted osmotic pumps available from ALZA, such as the ALZET osmotic pump.\nNasal delivery may be achieved by incorporation of the drug into bioadhesive particulate carriers ( less than 200 xcexcm) such as those comprising cellulose, polyacrylate or polycarbophil, in conjunction with suitable absorption enhancers such as phospholipids or acylcarnitines. Available systems include those developed by DanBiosys and Scios Nova.\nA noteworthy attribute of the compounds of the present invention in contrast to those of various peptidic compounds referenced in the background section of this application is the demonstrated oral activity of the present compounds. Certain compounds have shown oral bioavailability in various animal models of up to 90-98%. Oral delivery may be achieved by incorporation of the drug into tablets, coated tablets, dragxc3xa9es, hard and soft gelatine capsules, solutions, emulsions or suspensions. Oral delivery may also be achieved by incorporation of the drug into enteric coated capsules designed to release the drug into the colon where digestive protease activity is low. Examples include the OROS-CT/Osmet(trademark) and PULSINCAP(trademark) systems from ALZA and Scherer Drug Delivery Systems respectively. Other systems use azo-crosslinked polymers that are degraded by colon specific bacterial azoreductases, or pH sensitive polyacrylate polymers that are activated by the rise in pH at the colon. The above systems may be used in conjunction with a wide range of available absorption enhancers.\nRectal delivery may be achieved by incorporation of the drug into suppositories.\nThe compounds of this invention can be manufactured into the above listed formulations by the addition of various therapeutically inert inorganic or organic carriers well known to those skilled in the art. Examples of these include, but are not limited to, lactose, corn starch or derivatives thereof, talc, vegetable oils, waxes, fats, polyols such as polyethylene glycol, water, saccharose, alcohols, glycerin and the like. Various preservatives, emulsifiers, dispersants, flavorants, wetting agents, antioxidants, sweeteners, colorants, stabilizers, salts, buffers and the like are also added, as required to assist in the stabilization of the formulation or to assist in increasing bioavailability of the active ingredient(s) or to yield a formulation of acceptable flavor or odor in the case of oral dosing.\nThe amount of the pharmaceutical composition to be employed will depend on the recipient and the condition being treated. The requisite amount may be determined without undue experimentation by protocols known to those skilled in the art. Alternatively, the requisite amount may be calculated, based on a determination of the amount of target enzyme which must be inhibited in order to treat the condition.\nThe matrix metalloprotease inhibitors of the invention are useful not only for treatment of the physiological conditions discussed above, but are also useful in such activities as purification of metalloproteases and testing for matrix metalloprotease activity. Such activity testing can be both in vitro using natural or synthetic enzyme preparations or in vivo using, for example, animal models in which abnormal destructive enzyme levels are found spontaneously (use of genetically mutated or transgenic animals) or are induced by administration of exogenous agents or by surgery which disrupts joint stability.\nExperimental:\nGeneral Procedures:\nAll reactions were performed in flame-dried or oven-dried glassware under a positive pressure of argon and were stirred magnetically unless otherwise indicated. Sensitive liquids and solutions were transferred via syringe or cannula and were introduced into reaction vessels through rubber septa. Reaction product solutions were concentrated using a Buchi evaporator unless otherwise indicated.\nMaterials:\nCommercial grade reagents and solvents were used without further purification except that diethyl ether and tetrahydrofiran were usually distilled under argon from benzophenone ketyl, and methylene chloride was distilled under argon from calcium hydride. Many of the specialty organic or organometallic starting materials and reagents were obtained from Aldrich, 1001 West Saint Paul Avenue, Milwaukee, Wis. 53233. Solvents are often obtained from EM Science as distributed by VWR Scientific.\nChromatography:\nAnalytical thin-layer chromatography (TLC) was performed on Whatman(copyright) precoated glass-backed silica gel 60 A F-254 250 xcexcm plates. Visualization of spots was effected by one of the following techniques: (a) ultraviolet illumination, (b) exposure to iodine vapor, (c) immersion of the plate in a 10% solution of phosphomolybdic acid in ethanol followed by heating, and (d) immersion of the plate in a 3% solution of p-anisaldehyde in ethanol containing 0.5% concentrated sulfuric acid followed by heating, and e) immersion of the plate in a 5% solution of potassium permanganate in water containing 5% sodium carbonate followed by heating.\nColumn chromatography was performed using 230-400 mesh EM Science(copyright) silica gel.\nAnalytical high performance liquid chromatography (HPLC) was performed at 1 mL minxe2x88x921 on a 4.6xc3x97250 mm Microsorb(copyright) column monitored at 288 nm, and semi-preparative HPLC-was performed at 24 mL minxe2x88x921 on a 21.4xc3x97250 mm Microsorb(copyright) column monitored at 288 nm.\nInstrumentation:\nMelting points (mp) were determined with a Thomas-Hoover melting point apparatus and are uncorrected.\nProton (1H) nuclear magnetic resonance (NMR) spectra were measured with a General Electric GN-OMEGA 300 (300 MHz) spectrometer, and carbon thirteen (13C) NMR spectra were measured with a General Electric GN-OMEGA 300 (75 MHz) spectrometer. Most of the compounds synthesized in the experiments below were analyzed by NMR, and the spectra were consistent with the proposed structures in each case.\nMass spectral (MS) data were obtained on a Kratos Concept 1-H spectrometer by liquid-cesium secondary ion (LCIMS), an updated version of fast atom bombardment (FAB). Most of the compounds synthesized in the experiments below were analyzed by mass spectroscopy, and the spectra were consistent with the proposed structures in each case.\nGeneral Comments:\nFor multi-step procedures, sequential steps are noted by numbers. Variations within steps are noted by letters. Dashed lines in tabular data indicates point of attachment."}
{"text": "Additionally, in response to increasing social demands for energy and resource saving, there is an increasing need for a rubber material that is excellent in weather resistance and crack growth resistance to meet the requirements for improving durability of tires in order to reduce fuel consumption of automobiles. In addition, in view of the recent surge in price of butadiene, it is expected that the price of raw materials will even more dramatically rise in the future. As such, there is a need to make use of inexpensive olefin resource even in tire materials. To address this issue, some techniques have been proposed to improve weather resistance by blending an ethylene-propylene-diene rubber (EPDM) containing olefin as a principal component. These techniques, however, suffer from drawbacks, such as lower fracture resistance and lower crack growth resistance (see, for example, JP 2000-063639 A (PTL 1)).\nIn addition, applying a copolymer of a conjugated diene and a non-conjugated olefin to the compounded rubber results in fewer double bonds in the conjugated diene unit (conjugated diene compound-derived unit) in the copolymer as compared with conjugated polymers, and therefore, improved ozone resistance. In addition, the characteristics other than ozone resistance required when a rubber composition is applied to various applications (such as tires, conveyor belts or anti-vibration rubber) include good fracture resistance and good crack growth resistance. It is well known that coordination anionic polymerization using catalyst systems represented by a Ziegler-Natta catalyst allows for homopolymerization of olefins and dienes. However, it was difficult to provide efficient copolymerization of olefins and dienes using such polymerization systems.\nTo address this issue, for example, JP 2000-154210 A (PTL 2) discloses such a catalyst for polymerizing a conjugated diene that contains a transition metal compound of group IV of the periodic table having a cyclopentadiene ring structure, and also discloses α-olefins such as ethylene as monomers which can be copolymerized with the conjugated diene. However, PTL 2 does not provide a specific description of copolymerization of a conjugated diene compound and a non-conjugated olefin. Obviously, there is no description or suggestion that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nFor example, JP 2006-249442 A (PTL 3) discloses a catalyst for polymerizing olefins that consists of a transition metal compound such as a titanium compound and a co-catalyst, and also discloses a copolymer of an α-olefin and a conjugated diene compound. However, there is no description or suggestion in PTL 3 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 2006-503141 A (PTL 4) discloses a copolymer of ethylene and butadiene that is obtained by synthesizing ethylene and butadiene as a starting material by means of a special organic metal complex as a catalytic component. However, there is no description or suggestion in PTL 4 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 2000-086857 A (PTL 5) discloses a butadiene polymer having cis content of 92%, vinyl content of 6% and ethylene content of 3% or 9%. However, there is no description or suggestion in PTL 5 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 2000-154279 (PTL 6) discloses a rubber composition including: a butadiene-ethylene block copolymer having a cis content of 92% and ethylene segments in an amount of 4.8 mass % of the total; polybutadiene having a cis content of 95.2% and a vinyl content of 2.5%; and carbon black N220. However, there is no description or suggestion in PTL 6 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber.\nIn addition, JP 11-228743 (PTL 7) discloses an olefin-rich, olefin-diene copolymer. However, there is no description or suggestion in PTL 7 that a rubber which is excellent in weather resistance, fracture resistance and crack growth resistance can be obtained by mixing: a conjugated diene compound/non-conjugated olefin copolymer containing a conjugated diene compound-derived unit by at least 40 mol %; a conjugated diene-based polymer; and a non-conjugated diene compound/non-conjugated olefin copolymer containing an ethylene-propylene-diene rubber."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may or may not constitute prior art.\nA typical multiple speed transmission uses a combination of friction clutches or brakes, planetary gear arrangements and fixed interconnections to achieve a plurality of gear ratios. The number and physical arrangement of the planetary gear sets, generally, are dictated by packaging, cost and desired speed ratios.\nWhile current transmissions achieve their intended purpose, the need for new and improved transmission configurations which exhibit improved performance, especially from the standpoints of efficiency, responsiveness and smoothness and improved packaging, primarily reduced size and weight, is essentially constant. Accordingly, there is a need for an improved, cost-effective, compact multiple speed transmission."}
{"text": "It is common for molluscs such as oysters to be grown on a commercial scale in artificial environments in early stages of growth. Typically, the stock is then transferred to sea environments for what is known as “grow out”. Cultivation of oysters is accomplished by initially growing the lava in special tanks until they reach a juvenile size when they can be transferred to containers in intertidal or sub-tidal sites. Typically, the containers are in the form of baskets having perforated walls which enable water to flow through the baskets where the oysters use the water flow to obtain nutrients from the water as well as flushing impurities away.\nIn adjustable long line oyster cultivation, the baskets are typically suspended from lines that extend for hundreds of meters with multiple lines being attached in a parallel configuration according to an embodiment. These lines are typically made up of tensioned monofilament with a protective sheath or “drip tube” over the monofilament to reduce wear. Posts are used to support the tensioned line and are typically spaced at intervals along the line to allow for at least one basket to be suspended between adjacent posts. The drip tube typically extends the full length of the line. Special clips are used to attach the baskets to the drip tube, and can be attached both along the line (parallel arrangement) or across pairs of lines (perpendicular arrangement).\nIt has been recognised more recently that placing baskets in rougher waters can improve the growth of oysters. However, more aggressive water conditions have brought about problems not previously encountered in calmer water conditions. Most notably, it has been found that the clips used to rotatably attach the oyster baskets to the drip tubes are wearing, as are the drip tubes in areas where the clips attach. In rougher waters, it is not uncommon for clips and hence baskets to slide back and forth along the drip tube and collide with the support posts. Baskets have also been known to perform full 360 degree rotations about the line. This often results in winding of the drip tube and can also have adverse effects when, for example, winding occurs in opposite directions at opposite ends of the drip tube.\nFIG. 1 illustrates an example of a known clip 2 including at one end a means of rotatable attachment 4 to an oyster basket and at the other end a means of rotatable attachment 5 to a drip tube 6 which acts as a protective sheath for the line 7. The rotatable attachment 5 is simply a hook means with an opening of sufficient size to receive the tube when the hook is forced down on the tube and to maintain the tube therein. The result is that the clip is slideable along the tube, but not easily removable unless the clip is forced in a direction away from the line. It is the areas indicated by reference numeral 8 on the clip and reference numeral 9 on the drip tube which would typically wear and render the clips and tubes unusable over time. The attachment portion 5 of the clip is such that the drip tube 6 needs to be secured therein, but when this part is sufficiently worn, the tube 6 can easily become detached from the clip.\nThe skilled addressee would understand the expense and labour that is required to replace these clips, and moreso to have to remove clips and baskets suspended on a line to enable a full replacement of the line.\nIt is therefore an object of the present invention to alleviate some of the aforementioned problems, or at least provide the public with a useful alternative."}
{"text": "This invention relates to a personal digital cellular telecommunication system (PDC) using a code division multiple access (CDMA) system and, in particular, to a call control method for use in mobile communication to fix a radio zone upon call origination at a mobile terminal irrespective of a communication traffic and to a system therefor.\nIn a conventional mobile/hand-held cellular telecommunication system, a plurality of service areas (PDC/cell, PHS [personal handyphone system]/microcell) are continuously formed to enable connection of a communication line upon continuous movement of a mobile terminal. In each radio service area (may be referred to as a radio zone), a single radio base station (cell base station) accommodated in a mobile communication control station via a wired or a wireless line or channel (WILL system) is arranged. The mobile terminal moves between the radio zones from one to another. In a CDMA system for example, the radio base station notifies information such as an identifier assigned to the station and an available radio channel by continuously transmitting the information via a radio channel called a pilot channel.\nMoving in the radio zone of each radio base station, the mobile terminal at first receives the pilot channel and also receives and collects the information of peripheral radio base stations. Upon call origination, the mobile terminal originates a call to a particular radio base station having a strongest reception field strength (RSSI (Reception Signal Strength Indicator)) among the information of the radio base stations collected as mentioned above, and starts communication therewith. Each radio base station can use a plurality of radio channels and determines a radio channel to be used with reference to a radio wave propagation state between the station and the mobile terminal to communicate therewith.\nIf the radio wave propagation state of a current radio channel being used is deteriorated during communication, i.e., upon occurrence of decrease in reception field strength (RSSI) or deterioration in quality data (QI (Quality Indicator)) representative of presence/absence of word errors, the current radio channel is switched into another radio channel. Specifically, a mobile assisted handover procedure is executed to switch the current radio channel being involved in communication into another radio channel so that no line interruption occurs. Proposal is also made of a method (diversity handover) in which a same signal is transmitted from a plurality of radio base stations to a mobile terminal and the mobile terminal selects a particular reception signal among a plurality of reception signals (see Masuhito AKEBI et al xe2x80x9cHandover Control Schemes Using Micro-Diversity Techniques for Microcellular Communications Systemsxe2x80x9d, Technical Report of IEICE, RCS 93-86, 1994).\nIn a CDMA mobile communication system, diversity handover may be called soft handover. In the CDMA, mobile communication system of an IS-95 standard, transmission is carried out under closed/open loop control for fixing transmission power of the pilot channel. In addition, a particular mobile terminal controls connection of a radio base station upon soft handover by the use of a peripheral zone set including the information of the pilot channels of the peripheral radio zones and an active set of a combination of the radio base stations simultaneously connected to the mobile terminal.\nAs conventional mobile communication systems of the type, xe2x80x9cCDMA Radio Transmission System and Transmission Power Control Device and Transmission Power Controlling Measuring Device Used in the Systemxe2x80x9d is disclosed in Japanese Unexamined Patent Publication (JP-A) No. H10-56421 and xe2x80x9cMobile Communication Systemxe2x80x9d is disclosed in Japanese Patent No. 2739850.\nIf the control is carried out to fix the transmission power of the pilot channel in the conventional systems described above, variation in communication/control traffic in and out of each radio zone causes the fluctuation in reception quality of the pilot channel at the mobile terminal. This results in variation in area of the radio zone. In other words, it is required to provide a control method for fixing the area of the radio zone upon call origination at the mobile terminal irrespective of the communication/control traffic.\nIf the mobile terminal transmits an access request (or a connection request) to a particular radio base station for multiple access on the basis of the above-mentioned peripheral radio zone set, the particular radio base station as a requested station determines allowance/rejection of the multiple access request, taking into account the communication traffic at the station and the communication quality, and notifies the result of determination to the mobile terminal. However, production of such control signal also constitutes a factor of deterioration of the communication quality. In the above-mentioned call control method, a mobile terminal during soft handover and a mobile terminal originating a new call are similarly dealt with. This increase a call loss ratio.\nSystems disclosed as conventional systems in the above-mentioned publications have similar problems.\nIt is an object of this invention to provide a call control method in mobile communication and a system therefor, which are capable of starting stable communication.\nIt is another object of this invention to provide a call control method in mobile communication and a system therefor, which are capable of fixing a call origination allowable area irrespective of a downlink communication traffic.\nIt is still another object of this invention to provide a call control method in mobile communication and a system therefor, in which a mobile terminal can produce a multiple access request in correspondence to a propagation loss.\nIt is yet another object of this invention to provide a call control method in mobile communication and a system therefor, in which a radio base station can control whether or not the station acts as a multiple access destination depending upon the communication traffic of the station without carrying out useless call control.\nIt is a further object of this invention to provide a call control method in mobile communication and a system therefor, which is capable of assuring, in case where a call setup request is produced by a mobile terminal when a communication traffic at a radio base station exceeds a threshold level, a communication capacity for a requested call to decrease a call loss ratio by suppressing the possibility that the radio base station acts as a multiple access destination.\nAccording to this invention, there is provided a call control method in mobile communication in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations carries out:\na step of storing pilot channel transmission power level information into the pilot channel to be transmitted; and\na step of carrying out transmission by the pilot channel;\nthe mobile terminal carrying out:\na step of receiving the pilot channel from each of the radio base stations;\na step of calculating a difference between a reception field strength of the pilot channel of each of the radio base stations and the transmission power level information received;\na step of detecting from the difference a particular radio base station of a minimum propagation loss; and\na step of originating a call to the particular radio base station of a minimum propagation loss thus detected.\nAccording to this invention, there is provided a call control method in mobile communication in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations carries out:\na step of judging whether or not a downlink communication traffic is great or small;\na step of making the pilot channel to be transmitted store transmission power level information having a value greater than a pilot channel transmission power level actually transmitted in case where the communication traffic is small or making the pilot channel to be transmitted store the transmission power level information having a value smaller than the pilot channel transmission power level actually transmitted in case where the downlink communication traffic is great; and\na step of carrying out transmission by the pilot channel;\nthe mobile terminal carrying out:\na step of receiving the pilot channel from each of the radio base stations;\na step of calculating a difference between a reception field strength of the pilot channel of each of the radio base stations and the transmission power level information received;\na step of detecting from the difference a particular radio base station of a minimum propagation loss; and\na step of originating a call to the particular radio base station of a minimum propagation loss thus detected.\nAccording to this invention, there is provided a call control method in mobile communication in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations carries out:\na step of storing pilot channel transmission power level information into the pilot channel to be transmitted; and\na step of carrying out transmission by the pilot channel;\nthe mobile terminal carrying out:\na step of receiving the pilot channel from each of the radio base stations upon execution of soft handover when approaching its radio zone;\na step of calculating a difference between a reception field strength of the pilot channel and the transmission power level information; and\na step of determining a particular radio base station for multiple access with reference to the difference.\nAccording to this invention, there is provided a call control method in mobile communication in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations carries out:\na step of judging whether or not a downlink communication traffic is great or small;\na step of identifying the case where the downlink communication traffic exceeds a predetermined threshold level in the judging step;\na step of setting transmission power level information to a value greater than a transmission power level transmitted by the station in case where the threshold level is exceeded in the identifying step;\na step of storing the pilot channel transmission power level information into the pilot channel to be transmitted; and\na step of transmitting the pilot channel;\nthe mobile terminal carrying out:\na step of receiving the pilot channel from each of the radio base stations upon execution of soft handover when approaching its radio zone;\na step of calculating a difference between a reception field strength of the pilot channel and the transmission power level information; and\na step of determining a particular radio base station for multiple access with reference to the difference.\nAccording to this invention, there is provided a call control method in mobile communication in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations carries out:\na step of judging whether or not a downlink communication traffic is great or small;\na step of making the pilot channel store first transmission power level information having a value greater and smaller than a pilot channel transmission power level actually transmitted in case where the downlink communication traffic is decreased and increased in the judging step, respectively;\na step of making the pilot channel store second transmission power level information having a value greater than the transmission power level transmitted by the station in case where the downlink communication traffic is increased to exceed a predetermined threshold level and a call setup request is received from the mobile terminal; and\na step of transmitting the pilot channel;\nthe mobile terminal carrying out:\na step of receiving the pilot channel from each of the radio base stations;\na step of calculating a difference between the reception field strength of the pilot channel of each of the radio base stations and the first transmission power level information received;\na step of detecting from the difference a first particular radio base station of a minimum propagation loss;\na step of originating a call to the particular radio base station of a minimum propagation loss thus detected;\na step of calculating a difference between the reception field strength of the pilot channel and the second transmission power level information upon execution of soft handover; and\na step of determining a second particular radio base station for multiple access with reference to the difference.\nAccording to this invention, there is provided a mobile communication system in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations makes the pilot channel to be transmitted store pilot channel transmission power level information and transmits the pilot channel;\nthe mobile terminal receiving the pilot channel from each of the radio base stations, detecting a particular radio base station of a minimum propagation loss given by the difference between a reception field strength of the pilot channel and the transmission power level information received, and originating a call thereto.\nEach of the radio base stations comprises:\nstoring means for storing the pilot channel transmission power level information into the pilot channel to be transmitted; and\ntransmitting means for transmitting on the pilot channel the transmission power level information stored by the storing means;\nthe mobile terminal comprising:\nreceiving means for receiving the pilot channel from each of the radio base stations;\ndifference generating means for calculating a difference between the reception field strength of the pilot channel of each of the radio base stations that is received by the receiving means and the transmission power level information received;\ndetecting means for detecting a particular radio base station of a minimum propagation loss given by the difference calculated by the difference generating means; and\ncall originating means for originating a call to the particular radio base station of a minimum propagation loss that is detected by the detecting means.\nAccording to this invention, there is provided a mobile communication system in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations makes the pilot channel to be transmitted store transmission power level information having a value greater than a pilot channel transmission power level actually transmitted in case where a downlink communication traffic is small and transmits the pilot channel or makes the pilot channel to be transmitted store the transmission power level information having a value smaller than the pilot channel transmission power level actually transmitted in case where the downlink communication traffic is great and transmits the pilot channel;\nthe mobile station calculating a difference between a reception field strength of the pilot channel received from each of the radio base stations and the transmission power level information received, and originating a call to a particular base station of a minimum propagation loss detected from the difference.\nEach of the radio base stations comprises:\njudging means for judging whether the downlink communication traffic is great or small;\nstoring means for making the pilot channel to be transmitted store the transmission power level information having a value greater than the pilot channel transmission power level actually transmitted in case where the communication traffic is small in the judgement by the judging means or making the pilot channel to be transmitted store the transmission power level information having a value smaller than the pilot channel transmission power level actually transmitted in case where the downlink communication traffic is great; and\ntransmitting means for transmitting on the pilot channel the transmission power level information stored by the storing means;\nthe mobile terminal comprising:\nreceiving means for receiving the pilot channel from each of the radio base stations;\ndifference generating means for calculating a difference between the reception field strength of the pilot channel of each of the radio base stations that is received by the receiving means and the transmission power level information received;\ndetecting means for detecting a particular radio base station of a minimum propagation loss given by the difference calculated by the difference generating means; and\ncall originating means for originating a call to the particular base station of a minimum propagation loss that is detected by the detecting means.\nAccording to this invention, there is provided a mobile communication system in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations makes the pilot channel to be transmitted store pilot channel transmission power level information and transmits the pilot channel;\nthe mobile terminal receiving the pilot channel from each of the radio base stations upon execution of soft handover when approaching its radio zone, and carrying out multiple access with reference to the difference between a reception field strength of the pilot channel and the transmission power level information.\nEach of the radio base stations comprises:\nstoring means for storing the pilot channel transmission power level information into the pilot channel to be transmitted; and\ntransmitting means for transmitting on the pilot channel the transmission power level information stored by the storing means;\nthe mobile terminal comprising:\nreceiving means for receiving the pilot channel from each of the radio base stations upon execution of soft handover when approaching its radio zone;\ndifference generating means for calculating a difference between a reception field strength of the pilot channel received by the receiving means and the transmission power level information; and\ndetermining means for determining a particular base station for multiple access with reference to the difference generated by the difference generating means.\nAccording to this invention, there is provided a mobile communication system in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations makes the pilot channel to be transmitted store transmission power level information having a value greater than a transmission power level transmitted by the station in case where the downlink communication traffic exceeds a predetermined threshold level and transmits the pilot channel;\nthe mobile terminal carrying out multiple access with reference to the difference between the reception field strength of the pilot channel received from each of the radio base stations and the transmission power level information upon execution of soft handover when approaching its radio zone.\nEach of the radio base stations comprises:\njudging means for judging whether the downlink communication traffic is great or small;\nidentifying means for identifying the case where the downlink communication traffic exceeds a predetermined threshold level in the judgement by the judging means;\ngenerating means for generating the transmission power level information having a value greater than the transmission power level transmitted by the station in case where the threshold level is exceeded in the identification by the identifying means;\nstoring means for making the pilot channel store the transmission power level information generated by the generating means; and\ntransmitting means for transmitting on the pilot channel the transmission power level information stored by the storing means;\nthe mobile terminal comprising:\nreceiving means for receiving the pilot channel from each of the radio base stations upon execution of soft handover when approaching its radio zone;\ndifference generating means for calculating a difference between a reception field strength of the pilot channel received by the receiving means and the transmission power level information; and\ndetermining means for determining a particular base station for multiple access with reference to the difference generated by the difference generating means.\nAccording to this invention, there is provided a mobile communication system in which communication is carried out through a radio channel between a mobile terminal and each of a plurality of radio base stations accommodated in a mobile communication control station and each radio base station at first carries out notification of information through a pilot channel to the mobile terminal moving into a radio zone thereof, wherein:\neach of the radio base stations makes the pilot channel store first transmission power level information having a value greater and smaller than a pilot channel transmission power level being transmitted in case where the communication traffic is decreased and increased, respectively, and makes the pilot channel to be transmitted store second pilot channel transmission power level information having a value greater than the actual transmission power level transmitted by the station in case where the downlink communication traffic is increased to exceed a predetermined threshold level and a call setup request is received from the mobile terminal;\nthe mobile terminal originating a call to a first particular radio base station of a minimum propagation loss detected from a difference between a reception field strength of the pilot channel received from each of the radio base stations and the first transmission power level information received, and determining, upon execution of soft handover, a second particular radio base station for multiple access with reference to the difference between the reception field strength of the pilot channel and the second transmission power level information.\nEach of the radio base stations comprises:\njudging means for judging whether the downlink communication traffic is great or small;\nfirst storing means for making the pilot channel store first transmission power level information having a value greater and smaller than the pilot channel transmission power level actually transmitted in case where the downlink communication traffic is decreased and increased in the judgement in the judging means, respectively;\nsecond storing means for making the pilot channel to be transmitted store second pilot channel transmission power level information having a value greater than the actual transmission power level being transmitted by the station in case where the downlink communication traffic is increased to exceed a predetermined threshold level and a call setup request is received from the mobile terminal; and\ntransmitting means for transmitting the pilot channel in which the first or the second transmission power level information is stored by the first or the second storing means;\nthe mobile terminal comprising:\nreceiving means for receiving the pilot channel from each of the radio base stations;\nfirst difference generating means for calculating the difference between the reception field strength of the pilot channel of each of the radio base stations that is received by the receiving means and the first transmission power level information received;\ndetecting means for detecting the particular radio base station of a minimum propagation loss with reference to the difference calculated by the first difference generating means;\ncall originating means for originating a call to the particular radio base station of a minimum propagation loss that is detected by the detecting means;\nsecond difference generating means for calculating, upon execution of soft handover, the difference between the reception field strength of the pilot channel and the second transmission power level information; and\ndetermining means for determining the particular radio base station for multiple access with reference to the difference generated by the first or the second difference generating means.\nThe mobile communication system employs a code division multiple access system.\nIn the call control method in mobile communication and the system therefor according to this invention, the mobile terminal selects the radio base station of a minimum propagation loss represented by the difference between the transmission power level information from the radio base station and the reception field strength upon actual reception of the pilot channel, and originates a call thereto. As a result, it is possible to start stable communication.\nIn the call control method in mobile communication and the system therefor according to this invention, each radio base station transmits the pilot channel with the transmission power level information having a greater (smaller) value stored therein in case where the downlink communication traffic as measured at the station is small (great). The mobile terminal selects the radio base station of a minimum propagation loss represented by the difference between the transmission power level information from the radio base station and the reception field strength upon actual reception of the pilot channel, and originates a call thereto.\nAs a result, a call origination allowable area (service area) can be fixed irrespective of the downlink communication traffic. In other words, it is possible to eliminate useless call control processing and to prevent enlargement of the radio zone in case where the downlink communication traffic at a particular radio base station becomes smaller than a reference traffic and the signal quality of the pilot channel is relatively improved.\nIt is also possible to prevent contraction of the radio zone in case where the downlink communication traffic becomes greater than the reference traffic at the particular radio base station and the signal quality of the pilot channel is relatively deteriorated.\nIn the call control method in mobile communication and the system therefor according to this invention, the mobile terminal selects the radio base station of a minimum propagation loss represented by the difference between the transmission power level information from the radio base station and the reception field strength upon actual reception of the pilot channel, and produces a multiple access request. As a result, the mobile terminal can produce the multiple access request in correspondence to the propagation loss.\nIn the call control method in mobile communication and the system therefor according to this invention, the mobile terminal selects the radio base station of a minimum propagation loss represented by the difference between the transmission power level information having a value greater than the transmission power level and the reception field strength upon actual reception of the pilot channel in case where the communication traffic of the radio base station measured at the radio base station is greater than the threshold level, and produces a multiple access request thereto.\nAs a result, the radio base station controls whether or not it will be included in multiple access destinations depending upon the communication traffic without causing useless call control processing. Thus, it is possible to control the communication traffic in the station.\nIn the call control method in mobile communication and the system therefor according to this invention, the pilot channel with the first transmission power level information having a greater (smaller) value stored therein is transmitted in case where the communication traffic at the radio base station is small (great). Furthermore, in case where the communication traffic of the station is greater than a specific value and the call setup request is received from the mobile terminal, the pilot channel with the second transmission power level information having a greater value stored therein is transmitted. The mobile terminal selects the radio base station of a minimum propagation loss represented by the difference between the first transmission power level information and the reception field strength upon actual reception of the pilot channel, and originates a call thereto. Furthermore, the mobile terminal selects the radio base station of a minimum propagation loss represented by the difference between the second transmission power level information and the reception field strength upon actual reception of the pilot channel, and produces a multiple access request thereto.\nSpecifically, the radio base station controls whether or not it is included in multiple access destinations depending upon the communication traffic while the call origination allowable area is fixed without depending upon the downlink communication traffic, and controls the communication traffic at the station to impart the communication capacity which has been used for multiple access to the mobile terminal producing the call setup request.\nAs a result, it is possible to eliminate useless call control processing and to prevent enlargement of the radio zone in case where the downlink communication traffic at a particular radio base station becomes smaller than a reference traffic and the signal quality of the pilot channel is relatively improved. It is also possible to prevent contraction of the radio zone in case where the downlink communication traffic becomes greater than the reference traffic at the particular radio base station and the signal quality of the pilot channel is relatively deteriorated.\nThus, the radio base station controls whether or not it acts as the multiple access destination depending upon the communication traffic at the station without causing useless call control processing. In case where the communication traffic exceeds the threshold level and the call setup request is produced, the possibility of acting as the multiple access destination is decreased to assure the communication capacity for the requested call and to reduce the call loss ratio."}
{"text": "Field of the Invention\nThe present invention relates to an image acquisition apparatus including a fiber optic member and an imaging device, and an image acquisition method.\nRelated Background Art\nIn image acquisition by a solid-state imaging device including a plurality of pixels, a fixed pattern noise (FPN) caused by a sensitivity variation in each pixel due to a pixel structure, an imaging characteristic, and the like of the imaging device is generated in an acquired image. Such a fixed pattern noise becomes a cause of deterioration in an image acquired by the imaging device.\nAs a method of suppressing an influence of such a fixed pattern noise, a method called flat field correction (FFC) is proposed. The flat field correction is sensitivity correction for each pixel of an image acquired by an imaging device, and correction is performed in such a manner that an output from a pixel with low sensitivity becomes high and an output from a pixel with high sensitivity becomes low, whereby an image with sensitivity being uniform in the whole image is acquired (see, for example, Patent Document 1).    Patent Document 1: Japanese Patent Application Laid-Open Publication No. 2014-131619    Patent Document 2: Japanese Patent Application Laid-Open Publication No. H10-282243    Patent Document 3: Japanese Patent Application Laid-Open Publication No. 2003-21732    Patent Document 4: Japanese Patent Application Laid-Open Publication No. 2005-259176"}
{"text": "This invention relates to an electronic musical instrument of the type wherein such various characteristics of the musical tones generated as the color, volume and effect are controlled by the manipulation of a plurality of volume adjusting members, switches or the like mounted on the surface of a control panel, and more particularly an electronic musical instrument which can be readily played or performed by a novice player by making inoperative some one of a plurality of operating members, if desired, or by enabling a control based on preset information.\nIn a prior art electronic musical instruments, for example an organ, only tone levers (tone volumes) were provided to control tone color and footage for each keyboard but in a modern electronic musical instrument, as a result of developments in digital techniques such performances as automatic rhythm, automatic bass/chord, automatic arpeggio, etc. are added and switches, volume adjusting members or the like are also mounted on the surface of the control panel in addition to the tone levers. For this reason, the number of the operating members mounted on the panel surface has been increased to 100 or more so that such a high grade electronic musical instrument is difficult to play for an unskilled player."}
{"text": "1. Field of the Invention\nThe present invention relates to a patch for beautification having a skin-refreshing effect, a skin-clearing effect and/or a slenderizing effect.\n2. Description of the Related Art\nPlant extract is believed to have a skin-refreshing effect, a skin-clearing effect and/or a slenderizing effect. A lotion, gel or cream containing the plant extract has already been on the market. However, such conventional products have the following problems: for example, the effect is not clear; the effect is not sustained enough; the product is sticky when it is used; and clothes are dirtied by the product. Thus, a product that contains the active component(s) having the above-mentioned effects which acts on a desired site of the body effectively and sustainedly, and is handled with ease, has not yet been obtained."}
{"text": "Activity detection, both friendly and malicious, has long been a priority for computer network administrators. In known public and private computer networks, users employ devices such as desktop computers, laptop computers, tablets, smart phones, browsers, etc. to interact with others through computers and servers that are coupled to the network. Digital data, typically in the form of data packets, are passed along the network by interconnected network devices.\nUnfortunately, however, malicious activities can cause harm to the network's software or hardware, or its users. Malicious activities may include unauthorized access or subsequent unpermitted use of network resources and data. Network administrators seek to detect such activities, for example, by searching for patterns of behavior that are abnormal or otherwise vary from the expected use pattern of a particular entity, such as an organization or subset thereof, individual user, IP address, node or group of nodes in the network, etc.\nSecurity appliances are used in known systems to provide network security. The appliance approach involves installing security appliances (which are typically servers or computers configured for providing security) at one or more locations in the network. Once installed, the appliance monitors traffic that traverses the network. Functions provided by the appliance may include malware detection, intrusion detection, unauthorized access or unauthorized use of data, among others. Unfortunately, security appliances cannot easily be scaled to handle temporary or permanent increases in network traffic. Increased network traffic often requires a security vendor to perform an appliance swap or an equally time-consuming appliance upgrade. Appliances also tend to have only limited network visibility because they are typically configured to monitor data traversing the link on which a respective appliance is installed only. Such an appliance will be unaware of activities occurring on other network segments monitored by other appliances and thus cannot use the additional context information pertaining to activities occurring on other network segments to detect a cleverly-designed piece of malware that may be difficult to detect from purely localized information.\nInstalled software products, rather than security hardware appliances, provide another approach to security for data networks. These products, such as anti-virus or anti-malware software, typically are installed on terminal devices (e.g., desktop and laptop computers, tablets, or smart phones). Data traversing the network between the terminal device is monitored by the installed products to detect malware in either inbound or outbound data. Unfortunately, installed software products also do not perform well in terms of scalability or network visibility. Installed products tend to be disposed locally on the terminal devices and thus also tend to have fairly localized views of the data on the network. They also tend to be installed on hardware that cannot be upgraded easily."}
{"text": "The present invention relates to coated gall-resistant surgical saw blades. Orthopedic surgeons utilize saw blades of various shapes and configurations in the performance of surgical procedures. One such saw blade is known as a reciprocating saw blade and is designed to interact with a precision slot which guides the saw blade in reciprocating movements to cut bone tissue during the performance of surgery.\nCustomarily, a reciprocating surgical blade has a plurality of teeth on each of its two edges. These teeth are customarily provided in a set, that is, consecutive teeth are bent in alternate directions. As this set is provided on known surgical saw blades, the ends of the teeth extend upwardly and downwardly beyond the respective planes of the respective side faces of the blade. Thus, often, it is these ends of the teeth which guide the saw blade within the precision slot. Of course, inherently, this causes wear of the teeth quite prematurely and results in the surgeon having to change saw blades quite often during the surgical case.\nAnother type of saw blade is known as a sagittal saw blade and consists of a flat elongated blade with a proximal saw engaging slot and a distal transverse surface having a plurality of teeth thereon, with the teeth being flat in configuration, that is, not provided in a set.\nThese are but two examples of surgical saw blades which may form a part of the present invention.\nStudies have been conducted on the effect on bone tissue of blades which have been galled for whatever reason. In an article titled \"Orthopedic Saw Blades A Case Study\" by H. W. Wevers, et al, published in the journal of Arthroplasty, Volume 2 No.1, March 1987, this problem is discussed. The following is quoted from this article: \"Because these blades are used primarily for total knee arthroplasty, it is probable that the damage occurred from direct contact of the cutting edges with metal templates or instruments used in the operation. This type of damage had a direct influence on the mechanical work needed to operate the saw.\" Later in the publication, the following is stated: \"Excessive heat induces thermal damage to osteocytes and expands the zone of necrosis beyond that shown microscopically.\" Further, the following is stated: \"Smooth, accurately cut surfaces are recognized as an important factor for bone ingrowth into porous-coated prostheses. Such clean bone cuts enhance prosthetic fit and setting, therefore promoting an even load bearing to the bone, and improved alignment of the prostheses or osteotomies.\" Finally, the following is stated: \"Damage to blade cutting surfaces due to inadvertent contact with templates and instruments may be unavoidable with currently available techniques.\"\nA publication titled \"Avoiding thermal damage to bone: Machining principals applied to powered bone surgery, a literature review\", by Ray Umber, et al further discusses the problems attendant in the prior art. The following is disclosed therein: \"Thus, cutting with a dull tool, increases the amount of frictional heat generated, much of which is now located in the workpiece itself. With a dull tool not only is the surface of the workpiece increasing in temperature but also the cut is no longer clean.\"\nA further problem with prior art saw blades is also set forth in this publication. In particular, prior art saw blades are so designed that it is difficult to provide water to the site of the operation for cooling purposes and to remove bone chips which are generated during sawing. Due to present saw blade design, \"coolant can not reach the dissection site\". \"It is, therefore, important to allow chip relief and to allow a cooling fluid to reach the dissection site and the dissecting tool. Proper attention to the technique will allow chip relief and the introduction of cooling irrigation, resulting in healthy, living bone which will heal readily.\"\nIn a further aspect, applicants have recognized a need to provide a surgical saw blade with a hard coating bonded thereto which will protect the blade surface to thereby increase the lifespan of the blade.\nIn the past, applicants have attempted to coat a surgical saw blade with polytetraflurorethylene, better known by the trademark TEFLON. Experiments by applicants with TEFLON coated blades have revealed increased efficiency with marked reduction in galling. However, it was found that autoclaving as well as use of the blades in a precision slot would result in the coating stripping off and it was further discovered that chunks of coating would sometimes enter the surgical site. These chunks were very offensive to the surgical site since it was found that the TEFLON flakes were not bio-compatible with the tissues of the patients.\nApplicants also experimented with ceramic coatings such as calcium phosphate and aluminum oxide. Experiments with coating surgical saw blades with these materials failed for several reasons. Firstly, the melting point of the ceramics was higher than that of the blade material and, as such, during the coating process, the temper of the blade would be destroyed. Further, it was discovered that a blade coated with ceramic could not be bent or flexed without cracking the coating.\nThereafter, experiments were conducted with combinations of methyl methacrylate and ceramics. While these combinations were superior to the TEFLON, calcium phosphate and aluminum oxide, it was found that adherence to the metallic blades was inferior. Even the milling of grooves to hold the methyl methacrylate and ceramics to the blade did not solve the problems in a satisfactory manner.\nAs such it was concluded from these experiments that if the blades were to be coated to improve durability thereof and consequently prevent galling, coating with a metallic alloy would have to be employed.\nA need has consequently developed for a surgical saw blade which will not only be more durable in use, but which will reduce heat generation adjacent bone tissue while also allowing access of cooling and flushing water to the surgical site.\nThe following prior art is known to applicant:\n______________________________________ U.S. Pat. No. ______________________________________ D 30,478 to Earle 864,812 to Thullier 2,670,939 to Harp 3,517,670 to Speelman 4,036,236 to Rhodes, Jr. ______________________________________\nEarle discloses a grass cutting blade having a single rib extending outwardly from one face thereof. Of course, this is different from the teachings of the present invention even concerning the ribbed embodiment thereof in that the present invention contemplates ribs on opposed faces of a blade having teeth formed in a set with the ribs guiding the blade in a precision slot and with the blade coated with a metallic coating.\nThullier discloses a knife and other cutting blade in several embodiments. As the cross-sections demonstrate, each of the blades disclosed in this patent have differing cross-sectional thicknesses at different areas along the lengths thereof which would make it impossible to use these blades in conjunction with a precision slot. Furthermore, none of the embodiments of Thullier teach the use of a metallic coating bonded to the blade.\nHarp discloses a mixing paddle having a plurality of rib-like structures thereon which appear to be of differing thicknesses as best seen in FIG. 4. Thus, Harp is quite distinct from the teachings of the present invention as failing to disclose a cutting blade nor a coating therefor and as not being designed for use in conjunction with a precision slot, among other reasons.\nSpeelman discloses a blood-letting lancet having two longitudinal ribs extending from one face thereof and a point at one or both ends thereof. The present invention is distinct from the teachings of Speelman as including in one embodiment, ribs on opposed faces of an elongated blade having teeth on opposed sides thereof, which teeth are formed in a set configuration and the further provision of a coating bonded thereto.\nFinally, Rhodes teaches the concept of a surgical saw blade having teeth on one side thereof and a single elongated longitudinal rib extending outwardly from one face thereof. The lack of a plurality of ribs on opposed faces of the Rhodes, Jr. blade would inherently allow pivoting of the blade from side to side as it moves in a precision slot thus inherently damaging the tooth set.\nConcerning the prior art discussed above, as pertains to one embodiment of the present invention, while ribs are disclosed in these patents in conjunction with elongated blades, the ribs are disclosed only for strengthening purposes. While the ribs of one embodiment of the present invention inherently strengthen the elongated blade, their main purpose is to provide guidance of the blade within a precision slot. This is different from the teachings of the prior art."}
{"text": "Printed forms of margarine are generally made by an extrusion process wherein margarine, supercooled fluid form, is fed into a temperature controlled chamber under pressure for subsequent crystallization. Once crystallization or a \"hardening\" of the margarine has taken place, it can then be molded into the desired printed forms. The term \"printed forms\" or \"prints\" is a word of art that is used to define the shape of the margarine after it has been wrapped, and generally, refers to wrapped margarine having a rectangular bar shape.\nIn order to process margarine into print form, liquid margarine is taken from its mixing tank and pumped at a uniform rate to a heat exchanger where it is cooled to a paste-like consistency, and then to an expansion or compensating chamber or tube. From this expansion or compensating tube, the cooled margarine, which is still in a fluid state, is forced by the feed pump pressure into a retracted piston. Through the synchronous motion of a gear arrangement, the retracted piston and its cylinder are rotated 180.degree. and the piston is pushed forward to emit the collected margarine into a crystallization chamber. The piston is moved forward by new product being forced into the other side of the piston by means of the pressure head already established by the feed pump. The piston and cylinder are again rotated to emit the margarine into the crystallization chamber by the same process.\nThis step feeding process is continued until the crystallization chamber is filled with margarine. The margarine experiences a settling or crystallizing action which can be defined in terms of the time it takes a given amount of margarine to travel from the floating piston mechanism through the chamber and up to the molding or print chamber interface. It is during this time that the margarine crystallizes or solidifies and thus enables it to be molded into print forms. The crystallized margarine is forced into the mold chamber apparatus by the intermittent pressure built up in the chamber through the action of the floating piston in the step feeding process thereby causing the volume displacement of margarine into the mold.\nThe action of the floating piston is generally operated in synchronization with the molding operation, that is, the print mold chamber is in a receiving position to accept margarine when the piston is ejecting margarine into the crystallization chamber, hence the positive displacement of the margarine. Once the crystallized margarine is displaced into the mold, it is removed and forwarded to a synchronously operated wrapping apparatus and then onto a likewise synchronously operated packaging machine for assembling the wrapped prints into containers. This margarine processing operation is described in the brochure entitled \"Lynch Morpac Model `CT`\" put out by the Lynch Corporation of Anderson, Indiana.\nOne of the disadvantages of the above-identified system is that the functionality of the floating piston is severely limited, not only by its fixed dimensional size, but also because each piston must be externally adjusted for precise volume displacement of margarine into the mold. Hence, the volume displaced by each of the rotary floating pistons must be adjusted precisely to the volume of the mold chambers. Otherwise, excess margarine will be forced through the mold chamber vents, thereby reducing officiency and promoting maintenance problems due to leakage and subsequent margarine caking buildup outside of the system. Further complications exist in that sufficient margarine may not be available to fill the mold cavity thereby leading to the formation of incomplete or underweight prints.\nAnother disadvantage is that the 180.degree. rotation of the piston and its ejecting movement offers a mechanically complicated and expensive system that must be continually maintained. Leakages of the liquid margarine occur through the crevices of the piston and its chamber and eventually into the crevices formed by the chamber and housing thereby causing cleaning and sanitary problems.\nMore importantly, however, use of the free floating piston means for metering margarine offers the distinct disadvantage of limiting the length or capacity of the crystallization chamber because of the fixed mechanical drive arrangement operating from the molding apparatus main drive. Increasingly higher pressures are required to force the margarine emulsion into the piston chamber if higher operating speeds are desired, thus limiting the efficiency of the processing operation. Variations in crystallization chamber lengths caused by product and production demands are, therefore, restricted, especially when margarines are composed of the \"softer\" varieties of oils which require longer chamber lengths for crystallization."}
{"text": "In a modern retail store environment, there is a need to improve the customer experience and/or convenience for the customer. Whether shopping in a large format (big box) store or smaller format (neighborhood) store, customers often require assistance that employees of the store are not always able to provide. For example, particularly during peak hours, there may not be enough employees available to assist customers such that customer questions go unanswered. Additionally, due to high employee turnover rates, available employees may not be fully trained or have access to information to adequately support customers. Other routine tasks also are difficult to keep up with, particularly during peak hours. For example, shopping carts are left abandoned, aisles become messy, inventory is not displayed in the proper locations or is not even placed on the sales floor, shelf prices may not be properly set, and theft is hard to discourage. All of these issues can result in low customer satisfaction or reduced convenience to the customer. With increasing competition from non-traditional shopping mechanisms, such as online shopping provided by e-commerce merchants and alternative store formats, it can be important for “brick and mortar” retailers to focus on improving the overall customer experience and/or convenience.\nElements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions and/or relative positioning of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present teachings. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments of the present teachings. Certain actions and/or steps may be described or depicted in a particular order of occurrence while those skilled in the art will understand that such specificity with respect to sequence is not actually required. The terms and expressions used herein have the ordinary technical meaning as is accorded to such terms and expressions by persons skilled in the technical field as set forth above except where different specific meanings have otherwise been set forth herein."}
{"text": "1. Technical Field\nThe present invention relates generally to a cloth-like heating element with two pairs of electrical conductors and parallel circuits which are formed via a severing process.\n2. Related Art\nAn existing cloth-like heating element, Korean Utility Model Registration #1100067, is issued to the same applicant as the present invention, and its functions are as follows:\nAccording to predetermined design, a great number of resistance heating wires are arranged at some intervals between individual threads which wind up on the loom's beam, while a fine copper wire and a thread are used as the weft.\nThrough the weaving, the resistance heating wires and the copper wires are interlaced with each other to form the heating element, thereby the cloth-like heating element with parallel structure is completed.\nThe above mentioned cloth-like heating element, which has the benefits induced by parallel circuitry, is only in use with very low voltages (Dc12V, 24V), because in case of increasing the number of heating wires for the purposes of enhancing durability or heat balance of the heating element, the input voltage for the heating element must be relatively decreased.\nThis is in accordance with Ohm's law, which states that electrical resistance is inversely proportional to the number of heating wires connected in parallel.\nTherefore, the above mentioned heating element cannot function under high voltages (AC100V, 220V), which is a fatal weakness."}
{"text": "Conventional methods of refining lead bullion containing antimony have included the use of refining reagents such as air, oxygen-enriched air, lead oxides, nitre and caustic, and aluminum, at high temperatures in excess of 1100.degree. F. (600.degree. C.). In such methods, gaseous oxidizing reagents such as air and oxygen-enriched air are pumped into the agitated bullion where the oxygen present in these gaseous reagents reacts with the antimony to form antimony oxides. When using solid oxidizing reagents, such as lead oxides or nitre and caustic, the reagents are added while the bullion is vigorously stirred. At the high operating temperatures involved, the antimony is oxidized and forms a dross which rises to the surface. In such processes, a large amount of lead is also either oxidized along with the antimony to form a lead oxide dross or is physically entrapped in the dross which rises to the surface as a high lead-content dross.\nIn using aluminum for antimony removal, the refining process is carried out at a temperature of 1300.degree. F. The aluminum, upon addition, melts and is mixed into the bullion to react with the antimony to form an aluminum-antimony intermetallic which rises and is separated from the bullion as a dross. In this process, a considerable amount of lead is also physically entrapped in the dross.\nIn U.S. Pat. No. 1,786,908 to Hanak, alkali metals and caustic alkalies are used to separate antimony from lead and/or tin. He teaches that sodium will combine with antimony to remove any part of the antimony, by adding only that portion of sodium corresponding to the quantity of antimony to be removed. Hanak heats the metal or alloy to approximately 200.degree. C. (approximately 360.degree. F.) or more above its melting point and stirs a predetermined quantity of sodium into the metal. After mixing well, the temperature of the bath is somewhat lowered and sufficient caustic alkali is stirred in to dissolve the antimony-sodium alloy. The amount of caustic soda used is approximately four times the weight of the antimony extracted. The caustic sodium antimony alloy melt is run off, allowed to cool, and then treated by a water wash to dissolve the caustic soda and to react with sodium in the sodium-antimony compound to form the hydroxide with antimony being liberated in the metallic state.\nThus, in the Hanak process, the caustic soda added is four times the weight of the antimony extracted or about eight to ten times the weight of alkali metal added, and produces a molten caustic soda-sodium antimony slag which is difficult to handle with conventional batch refining equipment in a smelter and, in addition, the operating temperature of the process is 200.degree. C. or more above the melting point of the lead-antimony bullion, thereby increasing the required fuel cost and significantly decreasing the overall efficiency of his process.\nU.S. Pat. No. 3,607,232 to Lebedeff teaches a process for refining lead, particularly lead requiring softening and desilverizing, which comprises adding an alkali metal from the group consisting of sodium and potassium and mixtures thereof to a pool of the molten lead, agitating the molten pool to remove tellurium in the slag. The tellurium is removed from the molten lead at an extremely rapid rate at a temperature between the melting point of the lead and the temperature of the lead from the conventional pyrometallurgical lead softening process, i.e. at temperatures which ordinarily are below 1200.degree. F. (649.degree. C.) and which, depending on the softening process employed, may be as low as 850.degree. F. (454.degree. C.). The preferred operating temperature for carrying out the refining process is below 1200.degree. F. and above the melting point of the lead to be detellurized, preferably at a temperature above 650.degree. F. (343.degree. C.). The amount of sodium or potassium or mixtures thereof employed is equivalent to 0.3 to 0.8 pounds of sodium per pound of tellurium to be removed from the lead, with 1 pound of sodium being equivalent to 1.7 pounds of potassium. He also contemplates antimony removal during the refining process to remove most of the antimony due to the fact that where arsenic and antimony content of the lead is above about 0.02% by weight, relatively large quantities of antimony as well as arsenic are incorporated into the tellurium during the detellurizing procedure. The sodium or potassium or mixtures thereof are incorporated in the lead by adding them as a master alloy of lead which may be added in a liquid or solid form during agitation of the melt. The agitation required to reduce the tellurium in the lead to a given amount will vary inversely with the vigor of the agitation, i.e. with more vigorous agitation a short agitation period can be used and vice versa. The Lebedeff process is practiced at a temperature above 650.degree. F. and below 1200.degree. F. and most preferably in the range 800.degree. to 850.degree. F. for the removal of tellurium. Lebedeff also utilizes his process to soften bullion containing copper, bismuth, antimony, arsenic, selenium, silver and gold; however, the addition of sodium hydroxide and metallic sodium for the removal of tellurium from the bullion results in a bullion which retains the original percentage of copper, bismuth, arsenic, silver and gold although significant reductions are realized in the content of antimony and selenium retained in the bullion.\nThus, the Lebedeff process must be carefully controlled in order to remove the desired quantity of tellurium from the bullion and will also remove certain amounts of antimony and selenium, along with the tellurium; however, subsequent processing is required in order to further refine the bullion to remove copper, bismuth and other metal values.\nAll of the above prior art processes have shortcomings in that they are inefficient either in materials utilization, length of time to process the bullion, and/or energy consumption and fail to adequately remove other metal values in the initial processing, thereby necessitating extensive further softening procedures. In contrast, the process of the present invention refines the bullion within a controlled temperature range by the addition of sodium to remove antimony and substantial amounts of other common alloying metals and elements, such as arsenic, cadmium, copper, nickel, silver, sulfur, selenium, tellurium and zinc present in the lead bullion, which react with the sodium to form intermetallics and compounds which are recovered in the dross."}
{"text": "This invention relates to the precision measurement of angles by optical means. More particularly it relates to apparatus which utilizes an echelle or blazed grating in a spectrometer interacting with the optical output from a dispersive element on the object of interest to achieve great sensitivity of movement of the resulting patterns on a detector.\nThe ability to measure angles to submicroradian accuracy is necessary for a variety of advanced applications. Increased accuracy in the alignment of existing instruments will allow for increased accuracy in their measurements and actions. Ion beam mapping of microstructures could be made much more detailed. Ion beam etching of microelectronic arrays could produce smaller and faster devices. The need in space technologies is particularly acute. Long range space communications require very accurate alignment of antennae. The same is true if long range beam power transmission is to be feasible, e.g. microwave, laser or particle beam. This level of accuracy is also needed for specialized scientific missions such as gravitational wave detection and precision astronomy.\nAt present the state of the art is represented by the autocollimator. It is capable of just microradian accuracy, and the physics of the device restrict it past a certain level. What is needed is a low cost, compact device that can achieve even better accuracy in measuring angles than the autocollimator."}
{"text": "Numerical information can be displayed in a variety of ways, each of which is suited for displaying a particular type or quantity of data. Many options are available for presenting data, especially with the advent of word processors, spreadsheets, and graphics software programs for creating tables and graphs.\nIn order to communicate information decisions effectively to users, designers need to understand the tasks in which readers engage when they look at displays. Because the user's tasks require certain sensory, perceptual and cognitive operations, one has to consider what type of information the user is trying to ascertain in reading the graph. For example, a well-designed line graph makes it easier to discern the slopes of lines, whereas, if the user's goal is to determine actual values at certain points, differentiating slopes is of little value.\nThere are various factors in making a decision about how to present data. The first factor concerns the amount of data being presented. When presenting a small amount of data, authors should weigh the communicative benefits of tabular or graphical presentation against the reader's cognitive costs. Generally, in this case, tabular presentation is preferred. However, with a large amount of data, the reader's use of the data, especially the degree of precision the reader is likely to need, becomes important. When relations are more important than precise values, the data is probably best displayed in a graph.\nFinally, one must then choose the type of graph to present the information. This decision depends on both the characteristics of the readers and of independent and dependent variables. Generally, common graphs with which all readers are likely to have experience are used: line graphs, bar graphs, pie charts, and scatter plots. The choice of graph type depends on the readers' informational needs. Either a line graph or a bar graph is used if readers need to determine relative or absolute amounts. A line graph is used, for example, when a user needs to determine the rate of increase of a dependent (criterion) variable as a function of change in an independent (predictor) variable. A bar graph is used when a user needs to determine the difference between the means of the dependent variable across different levels of the independent variable. A pie chart is commonly used when a user needs to determine primarily proportions but not absolute amounts. However, the type of information each of these graph types can convey is limited. For example, a pie chart can only display one aspect of an underlying data set. If the underlying data set consists of two aspects, company profits as well as company losses, two pie charts are required. One for the company profits and one for the company losses.\nThus, given the above background, what is needed in the art are improved systems and methods for displaying data."}
{"text": "The discovery of the giant magnetoresistive (GMR) effect has led to the development of a number of spin-based electronic devices. The GMR effect is observed in certain thin-film devices that are made up of alternating ferromagnetic and nonmagnetic layers. In a typical device, the relative orientations of magnetic directions of the ferromagnetic layers define a binary state of the device. The resistance across a device is generally lowest when the magnetic directions of the ferromagnetic layers are in a parallel orientation and highest when the magnetic directions are in an antiparallel orientation. In a related device, the magnetic tunnel junction (MTJ), the nonmagnetic metal layer is replaced by a thin insulating layer. The tunneling current through this insulating layer can also be measured as a resistance that changes depending on the relative magnetization state of the two layers. The MTJ device generally produces a much larger resistance change than the GMR device.\nOne type of GMR device is commonly referred to as a “spin valve.” GMR devices, including spin valves, can be used as data storage elements in magnetic random access memory (MRAM) devices. In this regard, exemplary MRAM applications of GMR devices are described in U.S. Pat. Nos. 6,147,922; 6,175,525; 6,178,111; 6,493,258, and U.S. Pat. App. Pub. No. 2005/0226064, all of which are incorporated herein by reference.\nA spin valve or an MTJ device typically includes two or more ferromagnetic layers that are separated by a thin layer of a non-magnetic material (either a metal or insulator) and also includes an antiferromagnetic layer that “pins” the magnetization direction of one of the ferromagnetic layers. FIG. 1A illustrates (in a simplified form) the layers in a typical spin valve 10 as seen from a side view. As shown in FIG. 1A, the spin valve 10 includes ferromagnetic layers 12 and 14 separated by a nonmagnetic layer 16. In a typical arrangement, one of the magnetic layers is configured to be a fixed layer 14. The fixed layer 14 is adjacent to an anti-ferromagnetic layer 18, such that the magnetization direction of the fixed layer 14 is “pinned” in a particular orientation. The arrow in the fixed layer 14 indicates an exemplary pinned orientation, though, in general, the orientation could be pinned in either direction. Thus, the magnetization direction of the fixed layer 14 remains relatively fixed when operational magnetic fields are applied to spin valve 10. A second magnetic layer 12 is termed a free layer 12. In contrast with the fixed layer 14, the magnetization direction of the free layer 12 is free to switch between parallel and antiparallel orientations, as indicated by the double-arrow symbol in the free layer 12. By applying an appropriate magnetic field to the spin valve 10, the magnetization direction of the free layer 12 can be inverted while the magnetization direction of the fixed layer 14 remains the same.\nFIG. 1B shows a three-dimensional view of the spin valve 10 of FIG. 1A. As shown, the spin valve 10 has a hard-axis (short-axis) and an easy-axis (long-axis). In the absence of an applied magnetic field, the magnetization directions of both the free layer 12 and the fixed layer 14 run substantially parallel to the easy-axis.\nTypically, a magnetic memory, such as a magnetic random access memory (MRAM) or a magnetic read only memory (MROM), comprises memory cells that include a GMR or MTJ based memory bit. To determine a logical state of a particular memory cell within a magnetic memory, the magnetic memory will include a pair of pre-set reference elements that are used in comparison with a particular memory cell's memory bit. If the resistance measurement of the memory bit is below an average resistance of the pre-set reference elements, the memory cell is in a first binary state (e.g., a binary value of “0”). On the other hand, if the memory bit is above the average resistance, the memory cell is an opposite binary state, (e.g., a binary value of “1”). This manner of reading and writing a memory cell is generally referred to as a high/low scheme.\nIn general, the pre-set reference elements need to be set during manufacturing. If the reference elements are set incorrectly, or if the resistance of a particular memory cell's memory bit changes over time, a memory cell may fail to properly store data and thus become unusable."}
{"text": "1. Field of the Invention\nThe present invention relates to a finder device. More particularly, the present invention relates to the improvement of a finder device in a digital camera capable of displaying on a display screen both an image being viewed and an image being shot.\n2. Description of the Related Art\nCameras are commonly known, which form an optical image of a subject on an image pickup device from light coming from the subject and passing through a photographing optical system, photoelectrically convert the optical image into electric image signals through the image pickup device, and display the resulting image on the screen of an image display device, such as an LCD monitor, based on the electric image signals.\nAs an example of such a camera displaying a subject image, a camera is known, which is provided with an image pickup device for monitoring the subject field (image pickup device for live-view display) separately from an image pickup device for capturing the subject image. In this type of camera, a half mirror is arranged in a finder optical system to guide a light beam of the subject to the image pickup device for live-view display. In a live view mode that is one of the operation modes of the camera, the image pickup device for live-view display repeats imaging to display a series of live images on the screen of an image display device (for example, see Japanese Patent Laid-Open No. 2000-165730).\nA technique is also known which uses another image pickup device, different from an image pickup device for capturing a subject image, to monitor a finder screen (for example, see the U.S. Pat. No. 4,704,022).\nIt is common practice for conventional single lens reflex cameras (hereinafter simply called “SLR cameras”) to cause a superimpose display for displaying multiple focus points in a viewfinder or to display shooting information in the viewfinder using a display component for finder display.\nHowever, the camera described in Japanese Patent Laid-Open No. 2000-165730 is not to display shooting information and a subject image at a time on the screen of the display device.\nOn the other hand, when a finder screen is monitored through another image pickup device as described in the U.S. Pat. No. 4,704,022, there may be cases where the brightness of the display component for finder display does not fit the brightness of the subject on the finder screen."}
{"text": "1. Field of the Invention\nThis invention relates to a method for testing the quality of a silicon wafer. More particularly, this invention relates to a method for testing a semiconductor silicon single crystal bar grown from molten silicon by the Czochralski (CZ) method or the float zone (FZ) method to determine the presence or absence therein of a D defect which is one form of the point defect.\n2. Description of the Prior Art\nAs a measure to test the quality of a semiconductor silicon single crystal bar which has been grown by the Czochralski method or the float zone method, the procedure which comprises preparing a polished wafer (PW wafer) as a sample from the bar, oxidizing the surface of the silicon wafer with steam at a temperature in the range of from 1,100.degree. C. to 1,200.degree. C., cooling the oxidized silicon wafer to normal room temperature, stripping the oxidized silicon wafer of the oxide film with hydrofluoric acid, and selectively etching the wafer surface for about two minutes with a mixture of K.sub.2 Cr.sub.2 O.sub.7, hydrofluoric acid, and water has been prevailing. This method represents a substantially perfect simulation of the heat treatment to be performed in the formation of a device on the silicon wafer. Because of its ability to determine the presence or absence of an oxidation-induced stacking fault within the wafer during the formation of the device, this method has been heretofore accepted as an important testing technique.\nIt has been known that an aggregate of point defects are introduced into a silicon single crystal bar while the growth thereof is in progress [(1) Takao Abe, \"Applied Physics,\" 59 (1990), p. 272]. Among other defects, the A defect which is suspected to be an aggregate of interstitial silicon atoms and the D defect which is suspected to be an aggregate of vacancy type lattice defects are particularly popular. The aggregate of point defects is suspected to cause an oxidation-induced stacking fault (OSF) during the formation of the device. A need has been felt for a method capable of testing a silicon wafer to determine the presence or absence therein of the aggregate of point defects. The test of crystallinity which involves the heat treatment described above is problematic in that the test is impracticable unless the target of detection is an oxidation-induced stacking fault, that the test is incapable of directly detecting an aggregate of point defects to be formed during the growth of a crystal, that the test as a process spends a large amount of time, and that the apparatus required for the test is expensive. As a means of test for the detection of an aggregate of point defects, the copper decoration method has been adopted to date. This method comprises preparatorily applying an aqueous solution of nitric acid to the surface of a silicon wafer and then heating the coated silicon wafer at 1,000.degree. C. for several hours thereby decorating an aggregate of point defects with copper.\nThe test by the copper decoration method mentioned above for the detection of an aggregate of point defects is problematic in that the test necessitates a heat treatment, that the treatment has the possibility of altering lattice defects themselves and, therefore, constitutes itself a destructive testing and consumes a long time, and that the apparatus required for the test is expensive."}
{"text": "Field of the Invention\nThe disclosure relates to a fixing device to be used in an electrophotographic image forming apparatus, such as a copying machine, a laser beam printer, and a facsimile.\nDescription of the Related Art\nA fixing device with a film having a low heat capacity is known as a fixing device to be disposed in an electrophotographic image forming apparatus. Such a fixing device includes a cylindrical film, a sliding plate that slides while in contact with an inner surface of the film, a heater that heats the sliding plate with radiation heat, and a roller that nips the film with the sliding plate to form a nip portion, as discussed in Japanese Patent Application Laid-Open No. 2013-114057. In the nip portion, the fixing device heats a recording medium bearing an image while conveying the recording medium to fix the image on the recording medium. Thus, this fixing device has features of a short warm-up time and low energy consumption.\nJapanese Patent Application Laid-Open No. 2013-114057 discusses a sliding plate of a fixing device. The sliding plate has a curvature portion that is curved along an inner surface of a film on an upstream side of a nip portion in a recording medium conveyance direction. The curvature portion can preheat the film before the film enters the nip portion. Thus, the curvature portion has a merit that the film can be effectively heated.\nHowever, since a longitudinal end portion of the film extending to an outer side of a longitudinal end portion of a roller is not retained in the nip portion, an inner side of the film strongly slides against the curvature portion of the sliding plate. The inner side of the film is thus liable to be abraded."}
{"text": "The present invention is aimed at a symmetrical rolling bearing and in particular at a multi-row symmetrical rolling bearing.\nThe invention is described with reference to a wheel bearing unit for driven or non-driven wheels of motor vehicles, such as for example of pickup trucks, light trucks or SUVs (sports utility vehicles). It is however pointed out that the present invention may also be used in other rolling bearings.\nThe prior art discloses rolling bearings which have an outer ring, an inner ring and rolling bodies arranged between the rings. It is also known from the prior art to arrange the rolling bodies in a plurality of rows. For example, DE 100 60 638 A1 discloses a double-row rolling bearing. In this way, it is possible to distribute the loads acting on the bearing between a plurality of rows of rolling bodies, and therefore to a plurality of rolling bodies.\nThe use of balls in double-row rolling bearings as disclosed by DE 100 60 638 A1 is disadvantageous because this is unsuitable for high axle loads. The use of tapered rollers in double-row wheel bearings has the disadvantage that a high frictional torque is generated in this way.\nRolling bearings having more than two rows are relatively complex to produce and, in addition, it is necessary to determine geometric dimensions which are suitable for production.\nThe present invention is therefore based on the problem of providing a rolling bearing which is firstly suitable for high axle loads and which can secondly be produced and assembled with a justifiable amount of expenditure. In addition, it is intended to provide a rolling bearing which can be exchanged for a conventional rolling bearing with tapered rollers."}
{"text": "1. Technical Field\nThe present invention relates to a test apparatus and a test method.\n2. Related Art\nTo test electrical characteristics of a semiconductor device, such as an insulated gate bipolar transistor (IGBT), a metal-oxide semiconductor field effect transistor (MOSFET), and a diode, by applying a current to the semiconductor device, the current is conventionally applied to the semiconductor device while the device under test is connected to a high-speed current source as disclosed in, for example, Japanese Patent Application Publication No. 2000-241503.\nIt, however, becomes more difficult to implement a high-speed current source as the value of the current increases. A conventional test apparatus uses a special power source that can respond at high speed and apply a large current, but such a special power source is expensive.\nIn addition, even such a special power source requires a certain period of time until the current reaches a sufficient test current level. Some currents may flow into the device under test before the current of the current source reaches a sufficient test current level, resulting in raising the temperature of the device under test. If the temperature of the device under test exceeds a test condition, the reliability of the measured values may be compromised. Thus, it is required to realize a low-cost test apparatus that can response at higher speed and can apply a large current to a device under test."}
{"text": "Wooden components are, for example, constructed as plywood boards. In plywood boards, a number of thin wooden layers are glued and pressed together. The plywood boards have a relatively high strength to low weight ratio and can be easily processed, like cut or drilled. Another well-known board-like component is constructed as a so-called coreboard. Such coreboards have solid wood profiles which are glued together inside and which are covered by plywood boards on the outside. Such coreboards are generally used to construct solid furniture components or panelling for buildings.\nHowever, the well-known components made of wood and having generally flat shape and grains extending in one direction cannot meet all application requirements of high stability, low weight and high flexibility as well as requirements in savings of wood-like material and labor, when using them for instance in construction work.\nProposals has been made to improve the properties of the components made of wood, like plywood. For instance, a method of making a corrugated vegetable-fiber board from material which have sufficient plasticity when wet and becoming rigid upon being dried out is disclosed in U.S. Pat. No. 2,058,334 (Mason). In this prior art method the sheet having a thickness of about 3.2 mm is treated by several wetting, oiling, heating, drying and pressing stages. The method, however, includes the disadvantages in that the corrugation process of one single sheet is time consuming, requiring at least one day to produce and in most cases the producing time is even longer. Several heating stages at relatively high temperatures are required, which increases the need of energy and thus further the costs. Furthermore, since the corrugated panel is made from one flat panel by the above mentioned treatments, the corrugations are slightly retaining in use and especially under severe conditions of alternative wetting and drying out. A further disadvantage is that the resulting corrugated sheet is not thick enough to meet the requirements of wide variety of applications. Furthermore, the stiffness of the panels in direction perpendicular to the corrugations has also been found insufficient.\nGB Patent Application 2 087 793 discloses a corrugated building panel and method of manufacturing the same, said panel being formed of two plies of flexible paperboard or chipboard by steaming the two plies, applying an adhesive to facing side of one of the plies, corrugating the plies and thereafter immediately bonding the plies together. However, the method disclosed is capable to handle only sheets in the range of 0.9 to 1.3 mm in thickness and is merely related to paperboard sheets and not to sheets of wood-like fibrous material, like veneer. Furthermore, this proposal does not relate to the problem of stiffness of the produced panel or to the problem caused by the humid conditions, in which the panel might be used.\nIn GB Patent 1 360 105 a plywood construction is disclosed, in which plywood comprises a plurality of superposed interfitting corrugated veneers having their surfaces bonded together such that the grains of each veneer are crossing the ridges of the corrugations at an substantial angle to the ridges, the angle being always between 30 to 60.degree. and preferably 45.degree.. The method includes heat treatments of veneers with good bending properties in a steam filled atmosphere. Even though this solution has been able to give more rigid construction, there is a great loss of material when making corrugated plywoods, since the corners of the veneer sheets, which are substantially greater than the resulting corrugated panel, have to be cut off to achieve a corrugated panel having a square shape and corrugations extending from one edge to another. It is strongly pointed out in said GB Patent that if the angle of the grains is too great, i.e. over 60.degree., the veneers will splinter. An further disadvantage of that is that the appearance of the panel has nor fulfilled the aesthetic requirements due to the grains extending diagonally relative to the corrugations.\nThe object of the present invention is to overcome the disadvantages of the prior art proposals and to provide a new solution for a corrugated construction element and manufacture thereof, said element being sufficiently rigid for various applications and having good surface properties, but still being easy, quick and economical to manufacture and to process.\nAn another object of the invention is to provide a corrugated element comprising several and at least three layers of wood like sheets bonded to each other, said arrangements being such that at least one of the layers is arranged to be substantially perpendicularly placed relative to the corrugations.\nA still another object is that the layers are arranged such that grains of one wood sheet are essentially perpendicularly placed relative to adjacent sheets.\nA further object of the invention is to provide a corrugated panel of wooden sheets having an undulation of essentially small size.\nA further object of the present invention is to provide a corrugated construction element of wooden material, which keeps its shape as formed i.e. retains the shape of corrugations unchanged in use.\nA still further object of the present invention is to provide a corrugated component and a method for producing such a component such that the shape, size and spacing of the corrugations and the thickness of the layers of the element can be adjusted such that any desired form of undulations or corrugations can be provided.\nAn additional object of the present invention is to provide a method for continuous producing of the elements according to the invention.\nA further objective of this invention is to arrange the method in such a way that the component can be given an aesthetically attractive exterior design.\nA still further object is to provide a method by which the use of the prior art wetting or steaming of the boards is avoided."}
{"text": "1. Field of the Invention\nThe present invention relates to a data processing apparatus and method for performing multiply accumulate operations.\n2. Description of the Prior Art\nIt is common for data processing apparatus to be required to perform various floating point computations on data. To ensure a consistent approach in the way in which such floating point computations are handled by various data processing apparatus, a standard was produced in 1985 called the xe2x80x9cIEEE Standard for Binary Floating-Point Arithmeticxe2x80x9d, ANSI/IEEE Std 754-1985, The Institute of Electrical and Electronic Engineers, Inc., New York, 10017 (hereafter referred to as the IEEE 754-1985 standard). This standard defined, amongst other things, that a multiplication operation should finish with a rounding operation, and similarly that an add, or accumulate, operation should finish with a rounding operation. The IEEE 754-1985 standard further provided a definition of a number of rounding operations which would be considered to be compliant with the IEEE 754-1985 standard.\nIt has been found that general purpose processors are not well suited to the performance of floating point computations, and hence this has led to the development of specialised floating point units (FPUs) to handle such computations.\nOne particular floating point computation which is commonly required is a multiply-accumulate operation, whereby two numbers are multiplied together, and the product is then added to a third number. Multiply-accumulate operations were not specifically discussed in the IEEE 754-1985 standard, but rather multiplication and accumulate operations were discussed separately. Although a multiply-accumulate operation can be performed by executing a multiplication instruction followed by a separate accumulate instruction, such an approach is relatively slow.\nHence, there has been a great deal of interest in developing FPUs arranged specifically to perform multiply-accumulate operations with increased speed. An example of such a FPU is disclosed in U.S. Pat. No. 4,969,118, which describes a FPU developed by IBM to perform a multiply-accumulate operation. In accordance with the IBM technique, a partial multiplier produces a partial product of two numbers, and this partial product is then passed to adder circuitry for adding to a third number. Hence, the multiply-accumulate operation is xe2x80x98fusedxe2x80x99, in that the result of the multiplication is not independently determined prior to the accumulate operation. This approach significantly increases the speed of the multiply-accumulate operation.\nFurther, the multiplication is performed to an internal precision which contains all of the bits from the multiplication (for an nxc3x97n bit multiplication the result is 2 n bits) and the accumulation is then performed using all of the multiply bits. This provides a particularly accurate result, since no rounding is performed on the result of the multiplication before that result is used in the subsequent accumulation. However, it is apparent that this technique is not compliant with the IEEE 754-1985 standard since that standard defines that a rounding operation should be performed on the result of the multiplication.\nOther examples of FPUs designed specifically to increase the speed of a multiply-accumulate operation and/or reduce circuit complexity can be found in U.S. Pat. Nos. 5,241,493, 5,375,078, 5,530,663 and EP-A-0,645,699, U.S. Pat. Nos. 4,866,652 and 4,841,467. None of these documents are concerned with the issue of rounding, and in particular none are concerned with producing results which are compliant with the IEEE 754-1985 standard.\nAn alternative approach used in the MIPS R10000 product is to retain the multiplier and adder as separate logic units. When performing a multiply-accumulate operation, rounding is applied to the output of the multiplier unit, and this output is then input to the adder logic unit, with the result of the adder logic unit also being rounded. Whilst this enables an IEEE 754-1985 compliant result to be achieved for a multiply-accumulate operation, it does not retain the speed benefits to be obtained from a specialised logic unit arranged specifically to perform multiply-accumulate operations.\nIt is an object of the present invention to provide a data processing apparatus and method for efficiently performing a multiply-accumulate operation in response to a single instruction whilst producing a result which is equivalent to the execution of a separate multiply instruction incorporating rounding, followed by a separate add instruction incorporating rounding.\nAccordingly, the present invention provides a data processing apparatus for performing a multiply-accumulate operation A+(B*C) in response to a single instruction identifying said multiply-accumulate operation, comprising: a multiplier for multiplying values B and C to generate an unrounded multiplication result, the multiplier further being arranged to generate first data required for rounding determination; an adder for adding the unrounded multiplication result to a value A to generate an unrounded multiply-accumulate result, the adder further being arranged to generate second data required for rounding determination; determination logic for using the first and second data to determine one or more rounding values required to produce a final multiply-accumulate result equivalent to the execution of a separate multiply instruction incorporating rounding, followed by a separate add instruction incorporating rounding; and rounding logic for applying the one or more rounding values to generate the final multiply-accumulate result.\nIn accordance with the present invention, a multiplier is provided to generate an unrounded multiplication result by multiplying values B and C, and further to generate first data required for rounding determination. Similarly, an adder is provided to generate an unrounded multiply-accumulate result by adding a value A to the unrounded multiplication result, and further to generate second data required for rounding determination. Determination logic is then arranged to use the first and second data to determine one or more rounding values required to produce a final multiply-accumulate result equivalent to the execution of a separate multiply instruction incorporating rounding, followed by a separate add instruction incorporating rounding.\nBy this approach, dedicated multiply-accumulate logic can be provided to enable fast execution of a multiply-accumulate instruction, whilst producing a result which is compliant with the IEEE 754-1985 standard.\nIn one embodiment, the determination logic can be arranged to determine more than one rounding value to be applied at appropriate steps during the addition operation. However, in preferred embodiments, the determination logic is arranged to generate a single rounding value to be applied by the rounding logic to the unrounded multiply-accumulate result to generate the final multiply-accumulate result. This approach yields further speed improvements in the execution of the multiply-accumulate operation.\nIn preferred embodiments, the first data generated by the multiplier comprises guard and sticky bits, and the determination logic comprises first logic for determining a multiplier rounding value from the first data. Further, the first data preferably comprises one or more least significant bits of the multiplication result, which are also used in the generation of the multiplier rounding value. In such preferred embodiments, the determination logic further comprises second logic for determining the one or more rounding values from the multiplier rounding value and the second data.\nIn preferred embodiments, the adder comprises an alignment shifter for aligning the smaller of the value A and the multiplication result prior to performing the addition, and a detection unit for detecting whether the bits shifted out by the alignment shifter are all ones or all zeros. In this embodiment, the second data generated by the adder preferably comprises guard and round bits and the output of the detection unit.\nIt will be apparent that the multiplier could be arranged in a number of ways.\nHowever, in preferred embodiments, the multiplier comprises a multiplication unit for generating a partial result in carry-save format and a product adder for producing the multiplication result. Further, the multiplier preferably further comprises a sticky adder for adding together a number of least significant bits of the partial product generated by the multiplication unit, and a guard and sticky bit generator for generating guard and sticky bits from the output of the sticky adder.\nIt will be appreciated that there are a number of ways in which the final multiply-accumulate result may be generated using the one or more rounding values. In preferred embodiments, that values A, B and C consist of a mantissa and an exponent, which are manipulated separately. Hence, in preferred embodiments, the unrounded multiply-accumulate result comprises a mantissa and an exponent, the data processing apparatus further comprises an incrementer for incrementing the mantissa, and the rounding logic comprises a multiplexer for outputting either the mantissa or the incremented mantissa as a final mantissa of the final multiply-accumulate result, in dependence upon one of said one or more rounding values generated by the determination logic. In such preferred embodiments, the one rounding value preferably comprises a final increment signal input to the multiplexer to indicate whether the mantissa or the incremented mantissa should be output as the final mantissa.\nIt will be appreciated that when incrementing the mantissa, an overflow condition may arise. To account for this, the rounding logic preferably further comprises exponent incrementing logic for incrementing the exponent if the incrementing of the mantissa resulted in an overflow, and the incremented mantissa is selected as the final mantissa.\nIn preferred embodiments, the determination logic comprises a first converter for receiving the second data generated by the adder and compensating for any shift right normalization applied to the unrounded multiply-accumulate result by the adder. Further, the determination logic is preferably arranged to determine a multiplier rounding value from the first data and the determination logic preferably further comprises a second converter for applying the multiplier rounding value to the output of the first converter circuit. In such embodiments, the determination logic preferably further comprises a third converter for applying a predetermined rounding formula to the signals output from the second converter to generate the one or more rounding values.\nViewed from a second aspect, the present invention provides a method of performing a multiply-accumulate operation A+(B*C) in response to a single instruction identifying said multiply-accumulate operation, comprising: multiplying values B and C to generate an unrounded multiplication result, and generating first data required for rounding determination; adding the unrounded multiplication result to a value A to generate an unrounded multiply-accumulate result, and generating second data required for rounding determination; using the first and second data to determine one or more rounding values required to produce a final multiply-accumulate result equivalent to the execution of a separate multiply instruction incorporating rounding, followed by a separate add instruction incorporating rounding; and applying the one or more rounding values to generate the final multiply-accumulate result."}
{"text": "1. Field of the Invention\nThe present invention relates to a contaminating-element analyzing method using a total reflection X-ray fluorescence analysis and an apparatus of the same, in particular, to a contaminating-element analyzing method for allowing contaminating elements to be detected and the concentrations thereof to be precisely calculated and an apparatus of the same.\n2. Description of the Related Art\nThus far, as a method for nondestructively analyzing elements contained in samples, a fluorescent X-ray have been used. In addition, to improve the sensitivity of the analysis, total reflection X-ray fluorescence analyzing methods have been developed. The application of this method to contamination control for semiconductor production lines are being studied.\nAmong the total-reflection X-ray fluorescence analyzing methods, since energy dispersive-type fluorescent X-ray analyzing method allows a spectrum of a wide energy region to be measured, multi elements can be measured simultaneously with a single solid-state detector (SSD) disposed just above a sample. In addition, since the energy dispersive-type fluorescent X-ray analyzing method does not require a crystal monochromator, the SSD can be disposed at a position very close to a sample. Thus, the energy dispersive-type fluorescent X-ray analyzing method can provide higher sensitivity than for example wave dispersive-type fluorescent X-ray analyzing method.\nIn the conventional energy dispersive-type fluorescent X-ray analyzing method, if the concentration of a contaminating-element is high (for example, around 10.sup.11 atoms/cm.sup.2), the peak and escape peak of the element can be easily identified in the measured waveform of the fluorescent X-ray. An integrated intensity for obtaining the concentration of an element can be calculated with a low error.\nHowever, when a contaminating-element in the level ranging from 10.sup.8 to 10.sup.9 atoms/cm.sup.2 which is required in the semiconductor production lines is quantified by the energy dispersive-type fluorescent X-ray analyzing method, the following problems will take place.\n(1) Since the total count number is small, if the concentration of a contaminating element is low, the peak of the element in the measured waveform of a fluorescent X-ray disappears in the background. Thus, the identification of the peak of the element becomes very difficult.\n(2) Even if the peak of an element can be identified, it may be difficult to separate this peak from the escape peak of the element. Thus, when the integrated intensity is calculated, an error may take place. In addition, since the escape peak of an element also depends on the intensity of primary X-ray, the amount of contamination, and so forth, it is difficult to quantify the escape peak and feedback it.\n(3) In the semiconductor production lines, if a sample is Si substrate, it is difficult to identify the peaks of elements (such as Mg and Al) whose atomic numbers are close to the atomic number of Si regardless of the concentrations of their contaminations. In addition, sum peaks which are present at peak positions of n.times.Si--K.alpha. (where n is any integer) become interfering peaks.\n(4) To raise the sensitivity of analysis, an Primary X-ray which is provided with a rotating target is used so as to increase the output power. In this case, since the effect of Compton scattering of a primary X-ray is large, the peak of an element may disappear in the spectrum of the primary X-ray. For example, in W-L.beta..sub.1 of a Primary X-ray, which has a high excitation efficiency against a transition metal, the identification of the peak of Zn may be difficult.\nThe present invention is made from the abovedescribed point of view."}
{"text": "1. Field of the Invention\nAspects of the present invention relate to a metallic separator for a fuel cell, and more particularly, to a metallic separator for a fuel cell having excellent anti-corrosive properties and a low contact resistance.\n2. Description of the Related Art\nIn a fuel cell, a fuel such as hydrogen, natural gas, methanol, or the like is oxidized to produce electrons and hydrogen ions at an anode. The hydrogen ions produced in the anode pass through an electrolyte membrane to a cathode, and the electrons produced in the anode are supplied to an external circuit through a wire. The hydrogen ions that reach the cathode are combined with the electrons that reach the cathode through the external circuit and with oxygen supplied from the outside to form water.\nFuel cells have been promoted as next-generation energy conversion devices since they have a high electricity generation efficiency and are environmentally friendly. Fuel cells can be classified into polymer electrolyte membrane fuel cells (PEMFCs), phosphoric acid fuel cells (PAFCs), molten carbonate fuel cells (MCFCs), solid oxide fuel cells (SOFCs), etc., according to the type of electrolyte used. The operating temperature, materials of constitutional elements and the like vary according to the type of fuel cell.\nPEMFCs can be operated at relatively low operating temperatures, such as, for example, 80 to 120° C., and have a significantly high current density. Thus, PEMFCs can be used as power supplies for vehicles and homes.\nThe main constitutional elements of a PEMFC are a bipolar plate and a membrane electrode assembly (MEA). In order to make PEMFCs compact, light and economical, it would be desirable to improve particular aspects of PEMFCs such as the bipolar plate.\nThe MEA includes an anode in which the fuel is oxidized, a cathode in which an oxidizing agent is reduced, and an electrolyte membrane interposed between the anode and the cathode. The electrolyte membrane has ionic conductivity in order to deliver hydrogen ions generated in the anode to the cathode, and is an electric insulator in order to electronically insulate the anode from the cathode.\nTypically, the bipolar plate has channels through which fuel and air flow and functions as an electron conductor for transporting electrons between MEAs. Thus, the bipolar plate should be non-porous such that the fuel and the oxygen can be kept separated, and should have excellent electrical conductivity and sufficient thermal conductivity to control the temperature of the fuel cell. Furthermore, the bipolar plate should have sufficient mechanical strength to bear a clamping force applied to the fuel cell and should be corrosion-resistant with respect to hydrogen ions.\nConventionally, graphite has been the most common material used to form bipolar plates in PEMFCs, and channels through which fuel and air flow have typically been formed using a milling process. A graphite plate has sufficient electrical conductivity and resistance to corrosion to meet the requirements of a PEMFC. However, graphite plate and milling processes for shaping the plates are expensive. Further, graphite plates are typically brittle and thus, it is difficult to process bipolar plates with a thickness of less than 2-3 mm. Due to the difficulty in decreasing the thickness of bipolar plates made of graphite, it is difficult to decrease the size of a fuel cell stack consisting of several tens to several hundreds of unit cells.\nIn order to reduce the production costs and the thickness of bipolar plates, an attempt has been made to use metals as bipolar plate materials. Metals have most of the physical properties required for bipolar plates, and the costs of metals and processing thereof are reasonable. If metal can be used as the bipolar plate material, the costs of the bipolar plate can be reduced by 99% or more.\nHowever, a metallic bipolar plate may erode under the acidic conditions that are present inside a fuel cell, and thus, an oxidized film that imposes a high electrical resistance may be easily formed. As a result, serious problems such as membrane poisoning and increased contact resistance may occur. Corrosion of the metallic bipolar plate not only causes defects in the bipolar plate itself, but also poisons the catalyst and electrolyte due to the diffusion of metal ions into the electrolyte membrane. When the catalyst is poisoned, catalytic activity is decreased, and when the electrolyte is poisoned, the proton conductivity of the electrolyte is reduced, thereby resulting in deterioration of the performance of the fuel cell.\nIn addition, since the corroded metal dissolves, the contact between the separator and the MEA deteriorates and the electrical resistance is increased, resulting in deterioration of the performance of the fuel cell.\nThus, it has not yet been possible to use metallic bipolar plates in PEMFCs because of the likelihood of corrosion of the metal. For example, in a 1000-hour performance test, PEMFCs having bipolar plates formed of stainless steel, a Ti alloy, an Al alloy, and a Ni alloy, respectively, showed lower performance than a PEMFC having a graphite bipolar plate.\nTherefore, research into ways to improve the anti-corrosive properties of a metallic bipolar plate, such as a method of coating a surface with a material having an anti-corrosive property, has been carried out.\nFor example, a method of coating the surface of a bipolar plate composed of Ti or stainless steel with a material such as TiN that has excellent anti-corrosive properties and is electrically conductive is disclosed in Korean Laid-Open Patent Publication No. 2003-0053406.\nWhile an Al alloy or a Ti alloy is susceptible to forming an oxidized film, stainless steel, which is relatively less susceptible to forming an oxidized film and has high corrosion resistance, is an excellent alternative to graphite.\nThe above discussion with respect to the bipolar plate is also applicable to an end plate, a cooling plate, and a separator.\nAn end plate is an electrically conductive plate having channels for either a fuel or an oxidizing agent on only one side. An end plate having channels for a fuel is attached to an MEA disposed at one end of a fuel cell stack, and an end plate having channels for an oxidizing agent is attached to an MEA disposed at the other end of the fuel stack.\nA cooling plate is an electrically conductive plate having channels for a fuel or an oxidizing agent on one side and channels for a cooling fluid on the other side.\nA separator is commonly used when a flow field is formed in diffusion layers of an anode and a cathode and generally a bipolar plate is used when the flow field is not included. Advantageously, the separator may have low gas permeability, excellent electrical conductivity, and excellent anti-corrosive properties.\nHerein, the term “separator” is used generally to refer to the bipolar plate, the end plate, the cooling plate, and the specific separator discussed above.\nProblems relating to the bipolar plate or separator of a PEMFC have been described above, but such problems also occur in MCFCs, PAFCs, DMFCs, etc. Accordingly, the separator according to aspects of the present invention may be used in these types of fuel cells as well.\nBased on the above description, the development of a separator having improved anti-corrosive properties and contact resistance is desired."}
{"text": "The instant invention relates generally to devices for preventing the tampering of door locks and more specifically it relates to a protective guard for a lock cylinder in a door, which provides bolts that are concealed and adjustably connected on an inner side of the protective guard to mount the protective guard over the lock cylinder.\nThere are available various conventional devices for preventing the tampering of door locks which do not provide the novel improvements of the invention herein disclosed."}
{"text": "1. Field of the Invention\nThe present invention relates to amplification stages for providing a modulated supply voltage, and particularly to such amplification stages organised to provide higher power than can be achieved in a single stage by using an arrangement known as a ‘corporate structure’ amplifier.\n2. Description of the Related Art\nIt is known to provide a modulated power supply for providing a supply voltage to an amplification stage, such as a radio frequency (RF) amplification stage. An example of a particularly advantageous modulated power supply stage can be found in United Kingdom Patent No. 2398648.\nIn this advantageous modulated power supply stage there is provided an efficient technique for tracking the supply voltage to an RF amplifier in dependence on the RF input signal to be amplified by the amplifier. A first control loop tracks the envelope of the input signal, representing a desired supply voltage for the amplification stage, and selects one of a plurality of available supply voltages in dependence thereon. A second control loop tracks the envelope of the input signal and the actual output signal, and generates an error signal representing the difference there between. This error signal is combined with the selected supply voltage to provide an adjusted selected supply voltage for the amplification stage. The first control loop is a low frequency loop, and the second control loop is a high frequency loop.\nSuch an improved modulated supply stage offers significant advantages and improvements in delivering a highly accurate tracking (or tracked) power supply to an amplifier, with associated efficiency improvements.\nIt is also known, in amplifying signals such as RF signals, to split the input signal into portions, and then amplify each portion in a separate power amplifier, before recombining the amplified portions. The use of multiple power amplifiers to amplify a signal in this way enables a more powerful signal to be delivered than could otherwise be delivered. Typically such an amplified signal is delivered to an antenna of an RF transmitter.\nArrangements in which a signal is split in this way for amplification by multiple parallel amplifiers, before recombining, are known as ‘corporate structures’.\nIn the prior art highly accurate tracking power supply arrangement found in United Kingdom Patent No. 2398648, if the power supply stage is adapted for a corporate structure to deliver a power supply to n power amplifiers, then the power of each of the low frequency loop and the high frequency loop must be increased by a factor n. This has a negative effect on performance. The loops must be made physically larger. Layout problems may arise on the circuit board, with longer tracks needed to some power amplifiers. The maximum number of amplifiers which may be combined may also be restricted due to physical constraints of the power amplifiers on the feed network through which the modulated supply is delivered.\nAlternatively a single modulated power supply may be provided for each single amplifier stage. This requires a large amount of replication of circuits. This adds to the cost and size of the implementations, and thus introduces of its own inefficiencies. In this arrangement, to achieve the benefits of an efficient, highly accurate tracking power supply for each single amplifier requires a high component count and associated cost.\nAs noted above, in a corporate structure with n amplifiers and n associated modulated supplies, regardless of the specific modulated supply implementations, disadvantages of performance, cost and size may be encountered.\nIt is an aim of the invention to provide a technique to enable some or all of the benefits of high accuracy tracking power supplies to be obtained in arrangements in which multiple power amplifiers are implemented in parallel."}
{"text": "1. Field of the Invention\nThe present invention relates to a system and a method for providing hearing assistance to a user wherein audio signals from an audio signal source, which usually is a microphone arrangement, are transmitted by a transmission unit via a wireless audio link to a right ear unit and a left ear unit which are worn at or at least in part in the user's right ear and left ear, respectively, and which comprise means for stimulating the respective user's ear according to the transmitted audio signals.\n2. Description of Related Art\nUsually in such systems the wireless audio link is an FM (Frequency Modulation) radio link. The benefit of such systems is that sound captured by a remote microphone at the transmission unit can be presented at high sound pressure level to the hearing of the user wearing the ear units. In particular, the level of speech signals from the person using the transmission unit can be increased with regard to acoustic background noise.\nA typical application of wireless audio systems the receiver unit is connected to or integrated into a hearing instrument, such as a hearing aid. The benefit of such systems is that the microphone of the hearing instrument can be supplemented or replaced by the remote microphone which produces audio signals which are transmitted wirelessly to the FM receiver and thus to the hearing instrument. In particular, FM systems have been standard equipment for children with hearing loss in educational settings for many years. Their merit lies in the fact that a microphone placed a few inches from the mouth of a person speaking receives speech at a much higher level than one placed several feet away. This increase in speech level corresponds to an increase in signal-to-noise ratio (SNR) due to the direct wireless connection to the listener's amplification system. The resulting improvements of signal level and SNR in the listener's ear are recognized as the primary benefits of FM radio systems, as hearing-impaired individuals are at a significant disadvantage when processing signals with a poor acoustical SNR.\nMost FM systems in use today provide two or three different operating modes. The choices are to get the sound from: (1) the hearing instrument microphone alone, (2) the FM microphone alone, or (3) a combination of FM and hearing instrument microphones together.\nUsually, most of the time the FM system is used in mode (3), i.e. the FM plus hearing instrument combination (often labeled “FM+M” or “FM+ENV” mode). This operating mode allows the listener to perceive the speaker's voice from the remote microphone with a good SNR while the integrated hearing instrument microphone allows to listener to also hear environmental sounds. This allows the user/listener to hear and monitor his own voice, as well as voices of other people or environmental noise, as long as the loudness balance between the FM signal and the signal coming from the hearing instrument microphone is properly adjusted.\nAn example of an FM system is found in Canadian 2 422 449 A1 wherein the FM receiver unit is mechanically connected to a hearing instrument.\nSuch FM systems often are used in rooms. However, due to reflections in rooms the quality of the reception of the FM audio signals will vary depending upon head movement, position of the user in the room, positions and movement of other people or objects, etc. This varying quality manifests itself by a kind of a hissing noise and is especially noticeable in very small FM receivers as these receivers are built with very small antennas. These “holes” in the FM audio signal reception quality are an issue both with the current analogue FM systems as well as with the upcoming new digital systems.\nFurther, binaural hearing systems are already available, wherein there is provided a usually wireless link between the right ear hearing aid and the left ear hearing aid for exchanging data and audio signals between the hearing aids for improving binaural perception of sound. Examples of such binaural systems can be found in European Patent Application 1 651 005 A2, U.S. Patent Application Publication 2004/0037442 A1 and U.S. Pat. No. 6,549,633 B1. In European Patent Application 1 531 650 A2 a binaural system is described wherein in addition to the binaural link a wireless audio link to a remote microphone is provided. A similar system is described in WO 02/074011 A2.\nEuropean Patent Application 1 657 958 A1 relates to a communication system comprising a plurality of hearing aids between which audio signals may be exchanged via wireless links.\nU.S. Patent Application Publication 2005/0117764 A1 relates to a system comprising a right ear hearing aid and left ear hearing aid, each comprising a hearing coil for receiving audio signals from a telephone. The difference between the levels of the input signals of the two hearing devices at the hearing coils is measured and compared with a threshold value. If the difference in levels drops below or exceeds the threshold value, the respective hearing aid is switched to a telephone mode.\nU.S. Patent Application Publication 2004/0252852 A1 relates to a binaural acoustic beam forming system comprising a right ear hearing aid and a left ear haring aid wherein the voice-to-background noise ratio of the audio signal captured by the microphone of the right ear hearing aid and the audio signal captured by the microphone of the left ear hearing aid is determined and compared and wherein these audio signals are mixed prior to being supplied to the respective loudspeaker of the hearing aids, with the mixing ratio depending on the noise power ratio, i.e. the ratio of the voice to background noise ratios of the left ear hearing aid and the right ear hearing aid, respectively."}
{"text": "Substrates that include one or more layers of semiconductor material are used to form a wide variety of semiconductor structures and devices including, for example, integrated circuits (ICs) (e.g., logic processors and memory devices), radiation emitting devices (e.g., light emitting diodes (LEDs), resonant cavity light emitting diodes (RCLEDs), and vertical cavity surface emitting lasers (VCSELs)), radiation absorbing devices (e.g., optical sensors and solar cells) and switching/rectifying devices (e.g., power electronic devices). Such devices may comprise bipolar junction transistors (BJTs), power metal-oxide-field-effect-transistor (MOSFETS), thyristors, schottky diodes, junction field effect transistors (JFETs), insulated gate bipolar transistors (IGBTs) and PIN diodes. Such semiconductor devices are conventionally formed in a layer-by-layer manner (i.e., lithographically) on a semiconductor substrate.\nHistorically, a majority of such semiconductor substrates that have been used in the semiconductor device manufacturing industry have comprised thin discs or “wafers” of silicon material. Such wafers of silicon material are fabricated by first forming a large generally cylindrical silicon single crystal ingot, and subsequently slicing the single crystal ingot perpendicularly to its longitudinal axis to form a plurality of silicon wafers. Such silicon wafers may have diameters as large as about thirty centimeters (30 cm) or more (about twelve inches (12 in.) or more). Although silicon wafers generally have thicknesses of several hundred microns (e.g., about 700 microns) or more, only a very thin layer (e.g., less than about three hundred nanometers (300 nm)) of the semiconductor material on a major surface of the silicon wafer is actually used to form active devices on the silicon wafer.\nIt has been discovered that the speed and power efficiency of semiconductor devices can be improved by electrically insulating the portion of the semiconductor material that is actually used to form the semiconductor devices from the remaining bulk semiconductor material of the substrate. In addition, it has more generally been discovered that providing the semiconductor material used to form the semiconductor devices upon a base substrate, comprising one or more materials, allows for “engineering” of the properties of the semiconductor material.\nAs a result, so-called “engineered substrates” have been developed that may include a relatively thin layer of semiconductor material (e.g., a layer having a thickness of less than about three hundred microns (300 μm)) disposed upon one or more other material, such as, for example, dielectric material (e.g., silicon dioxide (SiO2), silicon nitride (Si3N4), silicon carbide (SiC), silicon (Si) or aluminum oxide (Al2O3)). Optionally, the layer of dielectric material may be relatively thin (e.g., too thin to enable handling by conventional semiconductor device manufacturing equipment), and the layer or layers of material over which the semiconductor material is disposed (i.e., the base substrate) may be thick enough to enable handling of the engineered substrate by manufacturing equipment.\nA wide variety of engineered substrates are known in the art and may include semiconductor materials such as, for example, silicon (Si), germanium (Ge), silicon carbide (SiC), III-V type semiconductor materials, and II-VI type semiconductor materials.\nFor example, an engineered substrate may include an epitaxial layer of III-V type semiconductor material on a surface of a base substrate, such as, for example aluminum oxide (Al2O3) (which is often referred to as “sapphire”). Using such an engineered substrate, additional layers of material may be foamed and processed (e.g., patterned) over the epitaxial layer of III-V type semiconductor material to form one or more semiconductor devices on the engineered substrate.\nWhen a layer of semiconductor material is epitaxially grown at elevated temperatures upon an engineered substrate, lattice strain may be induced in the crystal lattice of the layer of semiconductor material. The strain in the semiconductor material may result from a lattice parameter mismatch between the crystal lattice of the underlying materials of the engineered substrate and that of the semiconductor material being formed thereon (e.g., the underlying material has one or more lattice constants that differ from that of the semiconductor material).\nIn addition, a lattice strain may also be induced in the crystal lattice of the semiconductor material during epitaxial growth at elevated temperature due to a difference in the coefficients of thermal expansion (CTE) exhibited by the respective adjacent materials. For example, if the underlying engineered substrate has a mean CTE which is greater than that of the semiconductor material being grown upon the engineered substrate, the semiconductor material may be grown in a state of tensile strain. Such a state of tensile strain may increase during the growth of the semiconductor material as the layer thickness increases, and may eventually result in the formation of defects in the semiconductor layer. Such defects may include, for example, dislocations and cracks."}
{"text": "1. Technical Field of the Invention\nThe present invention relates to the use of encoded molecular sieve particles in an optical sensor analytical system.\n2. Description of the Prior Art\nThe use of optical fibers and optical fiber strands in combination with light absorbing dyes for chemical analytical determinations has undergone rapid development, particularly within the last two decades. Many of the recent improvements employing optical fiber sensors in both qualitative and quantitative analytical determinations concern the desirability of depositing and/or immobilizing various light absorbing dyes at the distal end of the optical fiber. In this manner, a variety of different optical fiber chemical sensors and methods have been reported for specific analytical determinations and applications such as pH measurement, oxygen detection, and carbon dioxide analyses.\nFiber optic sensors have been constructed that permit the use of multiple dyes with a single, discrete fiber optic bundle. U.S. Pat. Nos. 5,244,636 and 5,250,264 to Walt et al. disclose systems for affixing multiple, different dyes on the distal end of the bundle, the teachings of each of these patents being incorporated herein by this reference. The innovation of these patents is the placement of multiple chemical functionalities at the end of a single optical fiber bundle sensor. This configuration yields an analytic chemistry sensor that can be remotely monitored via the typically small bundle. The drawback, however, is the difficulty in applying the various chemistries associated with the chemical functionalities at the sensor\"\"s end; and, the functionalities are built on the sensor\"\"s end in a serial fashion. Not only is this a slow process, but in practice, only tens of functionalities can be applied.\nU.S. Pat. No. 6,023,540 by Walt et al. discloses a microsphere-based analytic chemistry system and method for making the same is disclosed in which microspheres or particles carrying bioactive agents are combined randomly or in ordered fashion and dispersed on a substrate to form an array while maintaining the ability to identify the location of bioactive agents and particles within the array using an optically interrogatable, optical signature encoding scheme. As a preferred embodiment, U.S. Pat. No. 6,023,540 teaches the use of a modified fiber optic bundle or array as a substrate to produce a high density array. The disclosed system and method have utility for detecting target analytes and screening large libraries of bioactive agents. The teachings of U.S. Pat. No. 6,023,540 are fully incorporated by reference herein.\nIn brief, the main limitation to present state-of-the-art technology whether it be through the use of microspheres, microbeads or particles, is the limited number of methods available to encode the array. Currently, polymeric-based microbeads are encoded by immobilized luminescent dyes only. In addition, there is a physical limitation to how many ultraviolet, visible, and near-infrared dyes can be used simultaneously to encode an array since the emission spectra of luminescence dyes are broad. Furthermore, present state-of-the-art technology utilizes only spherical microbeads. While an optional encoding avenue would be the use of microbeads with different diameters, this approach is limited by the difficulty in fabricating a large-scale batch of microbeads with a tight and uniform bead diameter distribution (which is the only way a plurality of spherical microbeads with different diameters could be employed in a reliable optical size-encoding scheme). In other words, present state-of-the-art microparticle-based analytical systems focus on the microparticle\"\"s chemical functionality and luminescent signature only. Therefore, while a small variety of silica-based and polymeric microspheres materials have been utilized, none of these microspheres offer size and shape selectivity.\nThe present invention represents an improvement over U.S. Pat. No. 6,023,540 as well as other comparable flow cytometric and fiber-optic sensor systems using microbeads, microspheres and/or microparticles. The key feature of the improvement is the added analytical performance features provided by shaped molecular sieve particles, namely optical encoding based on the molecular sieve particles\"\" macroscopic geometric shapes and increased selectivity based on the molecular sieves particles\"\" molecular-sized pore diameters i.e., pore sizes. Such encoded molecular sieve particles can provide at least a five-fold enhancement in tunable parameters for increasing the encoding possibilities of high throughput screening assays relative to the present dye-modified polymeric microsphere, microbead or microparticle standards.\nIt is an object of the present invention to provide a molecular sieve particle-based fiber-optic microwell array sensor, wherein a microwell array is etched onto an optical fiber bundle that is filled with molecular sieve particles having specific morphologies and pore sizes. The molecular sieve particles can be further modified with guest molecules including but not limited to, dyes, peptides, proteins, enzymes, antigens, antibodies, receptors, ligands, catalysts, nucleic acids, and oligonucleotides to form the basis of an optical chemical sensor or biosensor. These molecular sieve particle-based fiber-optic microwell array sensors also form the basis for combinatorial encoding and/or analysis. It is another object of the present invention to provide methods for synthesizing said molecular sieve particle-based fiber-optic microwell array sensors.\nA key feature of the invention is the added analytical performance features provided to the microwell sensors by the shaped molecular sieve particlesxe2x80x94namely optical encoding based on the molecular sieve particles\"\" macroscopic geometric shapes and increased selectivity based on the molecular sieves particles\"\" molecular-sized pore diameters. The term xe2x80x9cshaped molecular sieve particlesxe2x80x9d as used herein encompasses the macroscopic geometric shapes of the particles as well as their molecular-sized pore diameters. In addition to the ability to control and tune molecular sieve particle porosity, molecular sieve particles can be molecularly imprinted to further enhance selectivity to include the detection of analytes such as chiral/optically active molecules.\nOne key feature of the encoded molecular sieve particle-based optical sensor analytical system is the added optical encoding possibilities resulting from the plurality of the molecular sieve particle macroscopic geometric shapes that can be synthesized. In addition to spheres, molecular sieve particles can be synthesized with distinct gyroidal, discoidal, and hexagonal cylindrical shapes to increase the parameters by which optical shape-based encoding can be performed. Furthermore, the high surface area of the shaped molecular sieve particles enables encoding and detection in single microwells that is not possible with plastic beads or amphorous silica.\nPerhaps the most intriguing aspect of the molecular sieve particle-based fiber-optic microwell array sensor approach is the differences in the atomic compositions of molecular sieve materials in comparison to silica-based and polymeric microsphere materials. Specifically, silica-based and polymeric microspheres materials are comprised mainly of low atomic weight atoms such as carbon, hydrogen, nitrogen, oxygen, and silicon. Conversely, molecular sieve materials can be synthesized with a variety of atoms such as aluminum, titanium, iron, nickel, cobalt, germanium, gallium, boron, tin, selenium and other metals, metalloids, and non-metals. A variety of atoms can permit a number of new and alternative methods to be utilized for array encoding such as optical encoding by spectroscopic absorption techniques and/or energy dispersive and wavelength dispersive x-ray fluorescence techniques. Such approaches have the advantage of allowing the same limited number of luminescent dyes and the same number of macroscopic geometric shapes to be used multiple times, and thus, the total possible number of encoding combinations could be increased substantially.\nOne embodiment of the invention provides a chemical analysis method, comprising\npreparing separate subpopulations of shaped molecular sieve particles, each subpopulation carrying chemical functionalities that change optical signatures of said shaped molecular sieve particles in the presence of targeted analytes;\nencoding optical signature of the shaped molecular sieve particles in each subpopulation with a description of the chemical functionalities carried by that subpopulation;\ncombining the subpopulations to produce a system;\napplying the system;\ndetecting changes in the optical signatures indicative of the presence of the targeted analytes; and\ndecoding the optical signature of the shaped molecular sieve particles to identify the chemical functionalities.\nAnother embodiment of the invention provides an analytic chemistry sensor, comprising:\na bundle of optical fibers; and\na population of shaped molecular sieve particles carrying chemical functionalities at a terminal end of the fiber optic bundle. The terminal end may be the distal end or the proximal end of the optic fiber bundle.\nAnother embodiment of the invention provides the analytic chemistry sensor further comprising a source of electromagnetic radiation for exciting optically interrogatable chemicals bound to the chemical functionalities or for exciting optically interrogatable atoms that comprise the shaped molecular sieve particle.\nAn embodiment of the invention provides a method for constructing and using an analytic chemistry sensor, comprising:\nforming wells at terminal ends of optical fibers within a bundle;\ndistributing shaped molecular sieve particles carrying chemical functionalities within the wells; and;\nmonitoring a status of the chemical functionalities from an end face of the bundle.\nA Another embodiment of the invention provides a method for constructing and using an analytic chemistry sensor, comprising:\nforming wells at terminal ends of optical fibers within a bundle;\ndistributing shaped molecular sieve particles carrying chemical functionalities within the wells; and;\nmonitoring a status of the chemical functionalities from a proximal end of the bundle.\nAnother embodiment of the invention provides a composition comprising a plurality of optical fibers in an optical fiber array and a population of shaped molecular sieve particles, wherein said optical fibers have wells at a first terminal end of said fibers and a plurality of said wells contain at least one shaped molecular sieve particle.\nYet another embodiment of the invention provides a method of determining the presence of a target analyte in a sample comprising:\na) contacting said sample with a composition comprising:\ni) a substrate;\nii) a population of shaped molecular sieve particles comprising separate subpopulations, each subpopulation comprising:\n1) a chemical functionality for testing for interaction with a target analyte; and\n2) an encoding optical signature that can be used to identify said chemical functionality;\nwherein said shaped molecular sieve particles are distributed on said substrate; and\nb) determining the presence or absence of the target analyte.\nAn embodiment of the invention also provides a method of making a composition comprising:\na) forming wells at a terminal end of an optical fiber array; and\nb) distributing shaped molecular sieve particles within said wells, wherein said shaped molecular sieve particles comprise separate subpopulations, each subpopulation comprising:\ni) a chemical functionality for testing for interaction with a target analyte; and\nii) an encoding optical signature that can be used to identify said chemical functionality.\nAn embodiment of the invention further provides a sensor comprising:\na) an array of optical fibers;\nb) a population of shaped molecular sieve particles comprising separate subpopulations, each subpopulation comprising:\ni) a chemical functionality for testing for interaction with a target analyte; and\nii) an encoding optical signature that can be used to identify said chemical functionality;\nwherein said shaped molecular sieve particles are distributed on a first terminal end of said array; and\nc) a source of electromagnetic radiation.\nAn embodiment of the invention provides the optical fiber sensor further comprising a detector of electromagnetic radiation.\nAn embodiment of the invention provides a method for constructing and using an analytic chemistry sensor, comprising:\nforming wells at terminal ends of optical fibers within a bundle;\ndistributing shaped molecular sieve particles carrying chemical functionalities within the wells; and;\nmonitoring a status of a molecular sieve particle shape from an end face of the bundle."}
{"text": "A double patterning technology performed by plasma etching of CF4 gas has been known. The double patterning technology uses a wafer which includes a film to be processed, an organic film formed in a plurality of narrow lines on the film to be processed, and a Si oxide film that covers both the respective lines and the film to be processed which is exposed between the lines. In the double patterning technology, the Si oxide film is etched first to expose the respective lines of the organic film and the film to be processed. Then, the exposed organic film is selectively removed. Thereafter, the film to be processed is etched using the remaining Si oxide film as a mask. See, for example, Japanese Patent Application Laid-Open No. 2010-212415."}
{"text": "FIELD OF THE INVENTION\nThe invention relates to a method for controlling a drive train of a motor vehicle and to a drive train control system.\nIn conventional transmission control systems for motor vehicles, the transmission ratio of the transmission is set automatically as a function of the position of the accelerator pedal and of the speed of travel by using stored maps. Instead of using the position of the accelerator pedal as such, it is also possible to use a signal of a so-called electronic accelerator pedal (E-Gas for short) or some other variable derived from the engine torque demanded by the driver. In this case, various driving parameters and operating states of the motor vehicle are taken into account as is described in International Publication WO 97/01051. The respective gear to be selected or--in the case of continuously variable transmissions--the transmission ratio to be set is selected by control circuits operating by fuzzy logic methods. This logic describes expert knowledge gained from experience in the form of what is referred to as a basic set of rules. These rules are then used for regulating operations or control operations of the motor-vehicle transmission. The fuzzy logic circuit produces control signals, which, among other things, define the transmission ratio of the transmission and, in the case of automatic step-change transmissions, the gear to be selected.\nPublished, Non-Prosecuted German Patent Application DE 196 37 210 A1 discloses a conventional integrated drive train control system for a motor vehicle. The drive train control system interprets the position of the accelerator pedal and of the brake pedal as a wheel torque desired by the driver. It has a calculating device which receives the positions of the accelerator pedal and the brake pedal and which produces central control parameters for the drive unit and for those units of the drive train that exert a decelerating effect.\nThere are driving situations in which it would be expedient if actuation of the brakes by the driver were assisted by the drive train of the motor vehicle, when traveling downhill for example. This can, for example, be accomplished by the suppression of upshifts specified by a map, i.e. by retaining the existing transmission ratio.\nEuropean Patent Application EP 0 531 154 A2 discloses a transmission control system, with which the transmission ratio, which is determined as a function of the position of the accelerator pedal and of the speed of travel by using stored maps, can be corrected. For this purpose, a correction factor is determined by a fuzzy unit as a function of the driver's desire for braking, of the tractive resistance and of the current transmission ratio, and is added to the corresponding map value."}
{"text": "Aerofoils such as fan blades for a gas turbine engine are often manufactured as what is commonly known in the art as “wide chord blades”. The blades, typically formed of titanium, have a leading edge, a trailing edge and pressure and suction flanks extending between the leading and trailing edges and are manufactured from a series of plates which are laminated together before a gas injected between the sheets causes the blade to inflate. The blade is subsequently twisted and the external surfaces machined and polished to its finished form.\nThe flanks, leading and trailing edges may need to be machined and it is an object of the present invention to seek to provide an improved method of completing an operation on an aerofoil."}
{"text": "Motor vehicle restraint systems are known which include one or more load reducing webs. In some such restraint systems, a portion of an existing restraint web is looped upon itself and affixed via stitching or adhesive. When the restraint web is subject to a sufficient amount of force under load, the loops release, i.e., break apart, and the restraint web thus becomes longer and thereby transfers less load to the occupant than would otherwise be transferred with the shorter restraint web."}
{"text": "The present invention relates generally to an anti-pinch system.\nMechanisms for driving a closure, in particular for a vehicle, can be fitted with anti-pinch systems. These systems allow, for example, the direction of rotation of the drive motor to reverse when pinching is detected.\nThe drawback of these systems is that the inertia of the drive motor prevents instantaneous reversal of the window in the event of pinching. This can lead to an increase in pinching which exceeds the limits prescribed (for example, by Standard FMVSS 118) during the time required to reverse the direction of rotation of the motor, which prevents drive mechanisms from obtaining standard approval.\nDocument DE-A-26 15 511 describes a drive mechanism, in particular for a window regulator, including a clutch arranged between a reduction gear of a gear motor and a pinion for driving the window by an arm and sector mechanism. The clutch partially disengages the drive mechanism when an obstruction prevents the rise of the window, creating a resisting effort which opposes the torque of the gear motor. The clutch includes a spring designed to engage/disengage the mechanism at a set resisting torque threshold value. This device is simply a torque limiter. The mechanism maintains the pinching at the threshold value defined above. A need therefore exists for an anti-pinch system that is more effective and complies with existing standards."}
{"text": "Information drives business. Companies today rely to an unprecedented extent on online, frequently accessed, constantly changing data to run their businesses. Unplanned events that inhibit the availability of this data can seriously damage business operations. Additionally, any permanent data loss, from natural disaster or any other source, will likely have serious negative consequences for the continued viability of a business. Therefore, when disaster strikes, companies must be prepared to eliminate or minimize data loss, and recover quickly with useable data.\nCompanies have come to rely upon high-availability clusters to provide the most critical services and to store their most critical data. In general, there are different types of clusters, such as, for example, compute clusters, storage clusters, scalable clusters, and the like. High-availability clusters (also known as HA Clusters or Failover Clusters) are computer clusters that are implemented primarily for the purpose of providing high availability of services which the cluster provides. They operate by having redundant computers or nodes which are then used to provide service when system components fail. Normally, if a server with a particular application crashes, the application will be unavailable until someone fixes the crashed server. HA clustering remedies this situation by detecting hardware/software faults, and immediately restarting the application on another system without requiring administrative intervention, a process known as Failover. As part of this process, clustering software may configure the node before starting the application on it. For example, appropriate file systems may need to be imported and mounted, network hardware may have to be configured, and some supporting applications may need to be running as well.\nMany distributed computer systems utilize a centralized shared storage system for their provisioning. Thin provisioning is a mechanism that applies to large-scale centralized computer disk storage systems, SANs, and storage virtualization systems. Thin provisioning allows space to be easily allocated to servers, on a just-enough and just-in-time basis.\nThin Provisioning, in distributed computing system using a shared storage environment, is a method for optimizing utilization of available storage. It relies on on-demand allocation of blocks of data versus the traditional method of allocating all the blocks up front. This methodology eliminates almost all whitespace which helps avoid the poor utilization rates, often as low as 10%, that occur in the traditional storage allocation method where large pools of storage capacity are allocated to individual servers but remain unused (not written to). This traditional model is often called “fat” or “thick” provisioning.\nWith thin provisioning, storage capacity utilization efficiency can be automatically driven up towards 100% with very little administrative overhead. Organizations can purchase less storage capacity up front, defer storage capacity upgrades in line with actual business usage, and save the operating costs (electricity and floor space) associated with keeping unused disk capacity spinning.\nPrevious systems generally required large amounts of storage to be physically pre-allocated because of the complexity and impact of growing volume (LUN) space. Thin provisioning enables over-allocation or over-subscription.\nA volume manager is often used to manage large-scale centralized computer storage systems. However, problems exist where, in such systems, the thinly provisioned array and the volumes can have a very large number of file systems, where any volume can have any file system on it. Not being able to reclaim unused storage from the volume manager would require a customer to install/use specific services for each file system to be able to reclaim storage. Instead if the volume manager could reclaim unused storage of just about any file system, there would be just one method exposed by the volume manager."}
{"text": "The present invention relates to data storage systems, and more particularly, this invention relates to using performance data from tape heads and/or tape drives to predict future failure conditions.\nIn magnetic storage systems, magnetic transducers read data from and write data onto magnetic recording media. Data is written on the magnetic recording media by moving a magnetic recording transducer (Writer) to a position over the media where the data is to be stored. The Writer then generates a magnetic field, which encodes the data into the magnetic media. Data is read from the media by similarly positioning the magnetic read transducer (Reader) and then sensing the magnetic field of the magnetic media. Read and write operations may be independently synchronized with the movement of the media to ensure that the data can be read from and written to the desired location on the media.\nAn important and continuing goal in the data storage industry is that of increasing the density of data stored on a medium. For tape storage systems, that goal has led to increasing the track and linear bit density on recording tape, and decreasing the thickness of the magnetic tape medium. However, the development of small footprint, higher performance tape drive systems has created various problems in the design of a tape head assembly for use in such systems.\nIn a tape drive system, the drive moves the magnetic tape over the surface of the tape head at high speed. Usually the tape head is designed to minimize the spacing between the head and the tape. The spacing between the magnetic head and the magnetic tape is crucial and so goals in these systems are to have the recording gaps of the transducers, which are the source of the magnetic recording flux in near contact with the tape to effect writing sharp transitions, and to have the read elements in near contact with the tape to provide effective coupling of the magnetic field from the tape to the read elements."}
{"text": "This invention concerns improvements in printed circuit boards, and relates more especially to printed circuit boards of the kind bearing multiple conductor arrays such as may form data buses of a data processor, e.g. a microprocessor.\nWith the increasing miniaturisation of data processors there is a corresponding demand for printed circuit boards in which the area occupied by a given number of conductor tracks is reduced to a minimum. In the case of conductor arrays such as form data buses, and thus comprise a simple array of parallel conductor tracks, the number of tracks capable of being accommodated in a given area is theoretically limited only by the minimum width of conductive track giving acceptable resistance of the conductor and the minimum distance between adjacent tracks necessary to achieve acceptable isolation of the respective conductors. In practice, however, the termination of the conductor tracks at individual electronic components or electrical terminal connectors imposes limits upon the physical spacing of the conductors. Also, with high frequency and square wave signals such as occur in microprocessor systems, cross-talk between conductors increases significantly as the distance between conductors is reduced. Mitigation of these problems is therefore highly desirable in order to achieve improved standards of miniaturisation."}
{"text": "1. Field of the Invention\nThe present invention is related to detecting a flow of tissue and fluid aspirated from an eye. More specifically, the present invention is related to detecting a difference in viscosity between different materials being aspirated.\n2. Description of the Related Art\nDuring Ophthalmic Surgery both anterior and posterior surgery various materials, some transparent, are removed from an eye by aspiration, while other transparent materials are simultaneously infused into the eye. In anterior surgery the materials removed include a cataractous or clear lens, vitreous, cortical material and the infused material includes balanced-salt-solution (BSS) and viscoelastic. BSS and viscoelastic are infused to ensure that the eye stays inflated as an eye collapse during surgery can lead to catastrophic results. Posterior surgery material removed includes vitreous humor (vitreous), which is a clear gel filling the space between the lens and retina.\nIt is often difficult to distinguish between the various materials being removed because there is little visual difference between the different materials, except in the case of cataracts, which are easily identified. Also, a surgeon's ability to visualize the materials can be compromised by the eye's anatomy and inadequate illumination and magnification. The inability to distinguish between the materials being removed can lead to wasted time removing fluids infused into the eye, rather than removing the targeted natural tissue. This is especially true in posterior surgery during vitrectomy, where significant time can be spent removing BSS to ensure that all the vitreous is removed prior to operating on the diseased or torn part of the retina.\nOne difference between BSS and vitreous is the viscosity, where BSS is water and vitreous is a much more viscous gel.\nTherefore, it would be desirable if there were a feedback mechanism or alert that could inform a surgeon when the instrument removing material from the eye encounters a different material, particularly when a change in viscosity of the removed material is encountered."}
{"text": "The present invention relates to the general field of children's toys of the type having sensor activated behavior.\nThe typical type of toy with which the invention is specifically concerned is a doll that reacts to certain stimuli with an activated behavior, such as crying or making some other sound, waving a hand, or dancing. Although the description that follows is directed to a talking doll, the invention is not limited to this particular form of implementation.\nTo elicit the behavior from such dolls, an activation step is required. For instance, dolls having a pull string which, when pulled, causes a doll to speak simple preprogrammed words are well known. Additionally, toggle switches, buttons, loud noises and points on a doll responsive to pressure are also known to activate speech and motion in toys.\nLight can also trigger detectors to activate a toy or other device. For example, photosensors can detect darkness and, in response, turn on a light. Typically such devices monitor the amount of ambient light received by the photosensor, and when the amount of light received drops below a pre-established threshold, the light is turned on.\nSuch photosensors have been incorporated into dolls to elicit a response when the level of ambient light changes. For example, U.S. Pat. No. 5,501,627 describes a doll that has a photosensor located in each of the doll's eyes so that if either sensor is first exposed to ambient light, then deactivated by blocking the ambient light, such as by placing the doll's or the user's hand over at least one of the doll's eyes, a specific behavior, such as speech, is activated.\nWhile conventional sensors can be used to elicit behavior from a toy such as a doll, such toys do not respond in an identical matter in both light and dark conditions while also allowing the toy to respond to changes in ambient light unrelated to use of the doll. Since ease of use is a primary goal, a toy such as a doll whose behavior can be activated in both light and dark conditions is desirable."}
{"text": "Industrial commands for portable electronic devices including portable terminals have been increased as years go by, and portable electronic devices having various shapes and functions that coincide with consumers' requests have appeared.\nTypical portable terminals, such as portable phones, include display portions for displaying images or pictures thereon, and even portable terminals added with functions capable of viewing TV broadcasts through such display portions or visually obtaining various kinds of information using wired/wireless internet have appeared.\nFurther, various types of portable terminals, such as a folder type, slide type, and bar type, have been developed in accordance with consumers' requests and functional features of the portable terminals.\nHere, a folder type terminal means a terminal in which two bodies having liquid crystal screens formed thereon are hinge-operated to be unfolded in order to extend an available surface area, and has been continuously developed due to its superior structural features.\nRecently, with the development of thin film liquid crystal technology and miniaturization of other components, most portable terminals have generally been slimmed, and the demands for such slim portable terminals have been increased day by day.\nHowever, since the portable terminals in the related art have been developed on the basis of their external appearances that are more suitable to perform communication functions, such as voice call or message transmission, it is somewhat inconvenient to use them for multimedia services. For example, since the display devices of the general portable terminals are configured to be long in the vertical direction, it is inconvenient for users to view digital multimedia broadcasts or moving image files through such vertically long display devices.\nIn order to solve this problem, Korean Registered Patent No. 10-0909713, which is a registered patent granted to the inventors of the present invention, discloses “Hinge Device for Portable Terminal”. In an exemplary embodiment of this invention, in order to solve the above-described problem, the hinge device for a portable terminal is configured to rotate a cover of the portable terminal by a reference open angle and 360° in up-and-down direction while ultimately seeking the structural stability of the overall device.\nAccording to the invention of Korean Registered Patent No. 10-0909713, user convenience can be improved, and consumers' desires to use the display portion of the portable terminal in a relatively wide range can be satisfied. However, it still leaves room for improvement, and there is a need for the development of a new type hinge device and an electronic device using the same.\nFor example, in order to satisfy consumer needs for a large display screen even in a small-sized electronic device, it is required to provide an electronic device in which both two bodies that are hinge-rotated are provided with liquid crystal display screens formed thereon, and if the electronic device is unfolded, the two liquid crystal display screens are combined to form one large screen. Accordingly, it is also required to develop a hinge device that fits the required electronic device.\nThe inventor of the present invention has filed Korean Patent Application No. 10-2011-0047360 “Hinge Device and Portable Electronic Device Having the Same”, which has been granted to registration."}
{"text": "The present invention relates to the art of power driven threading machines and, more particularly, to a thread cutting die head having improved control for the receding of thread cutting dies in connection with cutting a tapered thread on a workpiece.\nPower driven thread cutting machines are of course well known and basically comprise a rotatable spindle and chuck assembly for rotating a workpiece to be threaded, and a thread cutting die head supported on a tool carriage for axial displacement relative to the workpiece to advance thread cutting dies supported by the die head into thread cutting engagement with the end of the workpiece. Generally, the die head assembly is pivotally supported on one side of the tool carriage for displacement between stored and use positions. In the use position, the axis of the die head is aligned with the workpiece axis, and the side of the die head opposite the pivotal mounting engages the corresponding side of the tool carriage to axially and vertically support the die head in the use position.\nIt is likewise well known in connection with cutting a tapered thread on a workpiece to release the thread cutting dies at the end of the thread cutting operation for displacement radially outwardly of the workpiece, or to provide for the thread cutting dies to progressively recede radially outwardly relative to the workpiece near the end of or throughout the thread cutting operation, followed by displacement away from the workpiece. In such die heads, the thread cutting dies are supported for radial displacement inwardly and outwardly relative to the die head axis and are so displaced by a cam plate coaxial with and pivotal about the die head axis. The cam plate is interengaged with the thread cutting dies in a manner whereby displacement of the cam plate in opposite directions about the die head axis displaces the thread cutting dies radially inwardly and outwardly thereof. The cam plate is spring biased to displace the thread cutting dies radially outwardly, and the thread cutting operation is initiated with the cam plate latched in a given angular position against the spring bias. Release of the cam plate provides for the spring bias against the cam plate to displace the latter to in turn displace the thread cutting dies radially outwardly of the workpiece. If such release of the cam plate is at the end of a thread cutting operation, the displacement of the cam plate and thus receding of the thread cutting dies, takes place immediately to retract the thread cutting dies to their outermost positions in which they are spaced from the workpiece. If such release is to provide for the progressive receding of the thread cutting dies near the end of or during the entire threading operation, angular displacement of the cam plate under the bias of the cam plate spring is controlled to achieve such progressive receding until the end of the thread cutting operation at which point the cam plate is immediately released for further angular displacement by the biasing spring to retract the thread cutting dies to their outermost positions away from the workpiece.\nHeretofore, progressive release of the cam plate near the end of a thread cutting operation, or during the entire threading operation has been achieved through the use of a lever actuated latch arrangement including a lever pivotally mounted on the die head and extending into the path of movement of a workpiece being threaded. The lever is engaged and displaced by the workpiece and in turn either disengages from the cam plate, immediately or progressively, or displaces another component which immediately releases or controls the progressive release of the cam plate to achieve gradual receding of the thread cutting dies near the end of or during the threading operation. Prior art arrangements for achieving progressive release of the cam plate in the foregoing manner near the end of the threading operation are shown, for example, in U.S. Pat. No. 4,288,181 to Sakaguchi et al, and Japanese Patent Application No. 57-20101. Prior art arrangements for achieving the progressive release of the cam plate in the foregoing manner during the entire threading operation are shown, for example, in U.S. Pat. No. 4,880,340 to Taki et al, Japanese Patent Laid-Open No. 57-66817 to Sakamoto, and Japanese UM Laid-Open No. 60-7931 to Gotov.\nIn the prior art arrangements providing for progressive release of the cam plate, the relative displacement between the component parts of the release mechanism is directionally transverse sliding displacement during which considerable force is applied between the relatively sliding surfaces resulting in rapid wear and, thus a loss in the accuracy and/or uniformity of threads being cut. More particularly in this respect, the lever or the component displaced thereby moves axially of the die head and is engaged by the cam plate or a component thereon which is biased by the cam plate biasing spring to pivot about the die head axis and thus transverse to the direction of movement of the lever or the component displaced thereby. Thus, during a threading operation, the relative displacement between the interengaging surfaces of the component parts is simultaneously axially and circumferentially with respect to the die head axis. It will be appreciated that such directionally transverse sliding interengagement between the component parts, which is repeated for each threading operation, results in rapid wearing of the component parts. Both the point of release of the cam plate at the end of a threading operation and control of the receding action of the thread cutting dies by the progressive release of the cam plate during a threading operation are important with respect to achieving accurate thread cutting and/or consistently uniform thread cutting. It will be appreciated that as the slidably engaging component parts of the release mechanism progressively wear, such accuracy and/or consistency is progressively reduced."}
{"text": "Patient data gathering devices used in emergency medical service (EMS) applications, for example, defibrillators, are typically complex and expensive. Thus, the number of devices employed in many areas is often relatively low. Nevertheless, the locations of such devices are typically tracked manually (for example, verbally or via an electronic log). Thus, it may be cumbersome to efficiently locate and redirect patient data gathering devices in certain situations, such as mass casualty situations.\nFurther, patient data gathering devices typically include wireless data transceivers to transmit medical information for storage in a remote medical database. However, some designs lack relatively long-range communication components (such as cellular or Wi-Fi transceivers) and in some situations, such as mass casualty situations, relatively long-range communication methods may be unavailable. In these cases, some patient data gathering devices turn to relatively short-range, device-to-device communication methods (such as near-field communications (NFC), Bluetooth, and the like) to transmit medical information. However, these devices may expend a significant amount of power by attempting to communicate with devices near the edge of their transmission range."}
{"text": "The invention relates to a device for attaching the severed parts of a band to one another. The trailing edge of a first terminating material band which is subjected to a treatment or to an operation in its downstream portion is connected to the leading edge of a second starting band. Included are a frame equipped with a means for the rotatable support of a supply roll or web for the first and second material bands and two cooperating, parallel pinch rollers. Each of the bands is guided along one of these rollers. A number of idler rollers are incorporated in the paths of the bands between each supply roll or web and the associated pinch roller. A device as described above comes in various embodiments, especially in relation to bands which are provided on both sides with an adhesive paste or adhesive tape.\nDepending upon the intended use of the material band (usually used for printing purposes), there exists a need to connect two successive band ends, e.g., by either overlapping or abutting their edges. This task imposes a high degree of care and precision on the machine operator who is called upon to carry out this task."}
{"text": "Iontophoresis is a method for introducing ionic substances into a body. The method utilizes direct electrical current to drive the ionized substances, such as drugs, through the intact skin or other body surface. This has proven to be very useful in numerous medical applications. U.S. Pat. Nos. 3,991,755 issued to Jack A. Vernon, et al and 4,141,359 issued to Stephen C. Jacobsen, et al disclose examples of iontophoretic devices and some applications of the devices. The iontophoretic process has been found to be useful in the administration of lidocaine hydrochloride, hydrocortisone derivatives, acetic acid, fluoride, penicillin, dexamethasone sodium phosphate and many other drugs.\nIn iontophoretic devices two electrodes are used. One electrode, called the active electrode, is the electrode at which the ionic substance is driven into the body. The other electrode, called the indifferent or ground electrode, serves to close the electrical circuit through the body. It will be appreciated by those skilled in the art that the active electrode must hold, contain or otherwise have available to it a source of the ionic substance. Thus, in the prior art the active electrode is generally relatively complex compared to the indifferent electrode.\nGenerally, prior iontophoretic drug delivery systems provide a single drug delivery rate. Such rate is obtained by applying a constant iontophoretic current designed to achieve a certain steady-state therapeutic concentration of drug in the body. With the use of such systems, there is a certain delay between the time that the iontophoretic maintenance current is initiated and when the desired therapeutic level of concentration is reached. Such delay may be, for example, thirty minutes from the time the iontophoretic current is initiated. In many cases, however, it is desirable or necessary that the iontophoretic drug reach therapeutic levels relatively fast. For example, where iontophoresis is used to deliver a narcotic pain killer, the patient often cannot tolerate a delay of even fifteen minutes. If the iontophoretic drive current is initially set at a relatively high level in order to encourage the rapid migration of iontophoretic drug into the bloodstream, the system will ultimately reach a steady-state level higher than desired or therapeutically safe. As a result, there is a need for an iontophoretic delivery system wherein therapeutic levels of drug concentration in the blood can be rapidly obtained while at the same time achieving a desirable steady-state maintenance level of administration."}
{"text": "The subject matter described herein relates generally to start-up methods and systems for start-up of a wind turbine.\nGenerally, a wind turbine includes a turbine that has a rotor that includes a rotatable hub assembly having multiple blades. The blades, which are also referred to as rotor blades, transform wind energy into a mechanical rotational torque that drives one or more generators via the rotor. The generators are sometimes, but not always, rotationally coupled to the rotor through a gearbox. The gearbox steps up the inherently low rotational speed of the rotor for the generator to efficiently convert the rotational mechanical energy to electrical energy, which is fed into a utility grid via at least one electrical connection. Gearless direct drive wind turbines also exist. The rotor, generator, gearbox and other components are typically mounted within a housing, or nacelle, that is positioned on top of a base that may be a truss or tubular tower.\nEach of the rotor blades is typically rotatably mounted on the hub using a respective pitch-drive system having a gearing supporting the rotor blade. Accordingly, the rotor blades can be adjusted relative to the wind direction. During partial load operation of the wind turbine, the rotor blades are typically in a fully powered position which may correspond to respective fully powered pitch angle values of 0° of the pitch-drive systems while pitch angle values of 90° may correspond to respective feathered positions. If the wind speed exceeds its rated value, the pitch angles of the rotor blades are typically varied so that the wind turbine can operate with fully rated power. On average, the pitch-drive systems typically operates the rotor blades in or close to the fully powered positions for more than about 50%, more than about 60% or even more than about 70% of the operating time. Accordingly, the fully powered positions of the rotor blades typically correspond to main operating positions in which and/or close to which the risk of gear wear of corresponding load portions of the pitch-drive systems is increased.\nThe load portions may be lubricated during a non-operating phase of the wind turbine, for example during a regular maintenance or in the absence of high enough wind. However, production phase may be long posing a risk of gear damage due to a non-surviving lubricating film. On the other hand, short maintenance intervals ensuring sufficient lubrication result in production losses and, thus, increase costs. This is of particular importance for coastal or offshore wind turbines.\nIt would therefore be desirable to provide systems and methods for improved lubrication of pitch drive systems."}
{"text": "1. Field of the Invention\nThe present invention relates to auditory devices, and more particularly to cable length adjustments for auditory devices.\n2. Related Art\nA certain category of hearing aid users are best assisted using a combination of hardware components for the ear which is being treated. A sound processor device, typically including one or more microphones, is provided in a behind the ear (BTE) device. A receiver unit is then provided in the ear canal, either partially or totally in the canal, in order to provide amplified acoustic stimulation for the user. The receiver is conventionally connected to the BTE device using an electrical cable which passes over or under the ear and then into the receiver.\nIt has been recognized that, for some users, the combination of a cochlear implant with acoustic stimulation in the same ear has significant benefits. In some such arrangements, a BTE device is provided, which is electrically connected via a cable to the receiver in the ear canal. The BTE device provides sound processing and outputs an electrical signal to the receiver via the cable.\nThe term ITE device is intended to encompass a device which is configured to fit wholly or partially within a recipient's ear, such as wholly or partially within the recipient's ear canal, whether for a cochlear implant system, hearing aid, bone conduction device, or other hearing prosthesis.\nAn electro-acoustic system is commercially provided under the Duet brand by Med-El. In this device, the receiver is integrated into the BTE device. The acoustic connection between the BTE device and the ear canal is a sound pipe for channeling the acoustic signals—the acoustic signals are generated within the BTE device.\nU.S. Pat. Nos. 6,748,094, 7,020,298 and 7,142,926, all assigned to ABC, describe the connection of the connection of the earhook to a BTE sound processor via different types of electro-mechanical connectors.\nOne issue with all such arrangements, in which a BTE device is electrically connected to a receiver in the ear, relates to the connecting cable. It is important to appreciate that the device is worn more or less constantly by the user and so issues such as comfort and cosmetic appeal are very important. The cable is the part of the system which is typically most vulnerable to damage. In many BTE hearing aid systems, an S shaped cable is used to provide some degree of adjustment and accommodation between the receiver and the BTE device. This allows for a small number of standard cable lengths, and earhook combinations, to be used for all users."}
{"text": "Active-matrix organic light emitting diode (AMOLED) displays are attracting attention due to several key advantages such as high efficiency, wide viewing angle, high contrast, and low fabrication cost. Among different technologies for implementation of AMOLED pixel circuits, hydrogenated amorphous silicon (a-Si:H) thin film transistor (TFT) is gathering more attention due to well established manufacturing infrastructure and low fabrication cost. However the threshold voltage (VT) of a-Si:H TFTs shifts over time with gate bias stress. If the current in the pixels depends on the VT of TFTs, VT shift causes degradation in the OLED luminance. This signifies the demand for pixel circuits and driving schemes that provide the OLED with a VT-independent current. Among different driving schemes, current programming has shown reasonable stability (A. Nathan et al., “Amorphous silicon thin film transistor circuit integration for organic LED displays on glass and plastic,” IEEE J. Solid-State Circuits, vol. 39, no. 9, September 2004, pp. 1477-1486). However, for small currents the programming time is large due to low field-effect mobility of a-Si:H TFTs and high parasitic capacitance of the data line. VT-compensating voltage-programmed pixels have smaller programming times (J. Goh et al., “A new a-Si:H thin-film transistor pixel circuit for active-matrix organic light-emitting diodes,” IEEE Electron Dev. Letts., vol. 24, no. 9, pp. 583-585, 2003) at the cost of imperfect compensation of VT.\nRecently, a driving scheme based on voltage feedback has been presented (S. Jafarabadiashtiani et al., “P-25: A New Driving Method for a-Si AMOLED Displays Based on Voltage Feedback,” Dig. of Tech. Papers, SID Int. Symp., Boston, pp. 316-319, May 27, 2005). The method provides proven stability and faster programming than the current-programming scheme. However, it is not fast enough to fulfill the demands for high-resolution large displays.\nIt is therefore desirable to provide a method and system that enhance the programming speed of a light emitting device display."}
{"text": "1. Field of the Invention\nThe present invention relates to an information processing system, and particularly to an information processing system for use in the design of a lens system.\n2. Description of the Related Art\nWhen an optical system is simulated by using an optical CAD (computer aided design) system, the data regarding the optical lens system should be defined in the CAD system as a matter of course.\nThe data with respect to the optical lens system that should be defined includes arrangement data such as the radius of curvature of a lens, the distance between lens faces, and the refractivity, which depends on the lens material simulated, and also includes evaluating position data such as the position of an object, and the evaluating image position.\nWith recent expanding applications of the optical technologies, various types of optical lens systems have been developed.\nIn the optical CAD system used for this purpose, it is required that various kinds of data associated with an optical lens system be easily defined in a manner which can fit with the actual development approach.\nTaking a zoom lens as an example, a method of data definition with a conventional CAD system will be described below.\nOne known CAD system for such a purpose is the CAD system CODE-V, which is available from Optical Research Associates.\nIn this optical CAD system, the data associated with a zoom lens may be defined as follows:\n(1) Input information on the evaluating light such as the wavelength, F-number, angle of view, and ratio of pupil diameter; PA0 (2) Input data representing the lens configuration such as the radius of curvature, the distance between lens faces, lens material, as well as the position and diameter of the diaphragm, which are associated with the respective lenses comprising the optical lens system; and PA0 (3) Evaluate positions such as the position of an object, and the evaluating position of an image.\nIn the case of an optical lens system of a camera, the position of an object and the evaluating position of an image correspond to the position of an object to be photographed and the position of a film plane, respectively.\nThen, the data of the optical lens system that has been defined in such a way as described above is used as the reference data to represent the standard conditions for describing a plurality of conditions, i.e., the zooming condition in this case, of the optical system. In this definition, the lens faces whose positions are variable are specified by their numbers, and variable parameters are used as the face distances belonging to respective lens faces.\nIn this optical CAD system, data regarding the evaluating positions such as the position of an object and the evaluating position of the image can be defined such that only one common set of the data need be given for all optical lens system conditions.\nIn another known optical CAD system called OPTEDGE, which is available from Minolta Camera Co., Ltd., data regarding evaluating positions such as the position of an object and the evaluating position of the image can be defined such that one set of the data may be given separately for each optical lens system condition.\nIn this system, as in the CODE-V system, initial input data regarding the optical lens system is used as the reference data representing the standard conditions, and based on this reference data, a plurality of conditions, i.e., the zooming condition, of the optical system can be defined in such a way that the lens faces whose positions are variable are specified by their numbers, and the variable parameters are used as the face distances belonging to respective lens faces.\nThe position information of the optical system represents the vertex position and the inclination of each lens face of respective optical elements which comprise the complete optical system.\nWhen an optical system that is decentered is designed using a conventional simulation system, the arrangements of the optical elements should be given for each lens face by defining its 3-dimensional location and its inclination and the like.\nReferring to FIG. 47, one example will be described below.\nIn the optical system shown in FIG. 47, the coordinate of a face vertex (solid dots along the axis in FIG. 38) should be given for each face to define the position data of each optical element.\nIn this figure, the origin of an absolute coordinate system is assigned to (0,0,0) corresponding to the face vertex of lens face 1. However, the origin can be assigned to an arbitrary position.\nMore specifically, each arrangement of lens faces 1 through 13 is described by using 3-dimensional coordinates in the absolute coordinate system (the coordinates seen from the face vertex of the face 1). It will be a very troublesome task for a user to input the data after calculating the 3-dimensional coordinates for all faces ((x2, y2, z2), . . . , (x13, y13, z13)).\nFurthermore, as various kinds of optical systems have appeared, decentered lens systems become used more often, and it becomes desirable for a simulation system to include 3-dimensional information.\nWhen a decentered optical system is designed with a conventional simulation system, the origin, the inclination of the coordinates, and the rotation data should be given for each lens face.\nIn conventional techniques, there are no simulation systems in which a plurality of lens faces are combined into one unit so that the unit has information representing the decentering.\nYoshinari Matsui describes in his book entitled \"The Design of a Lens System\" that the development of a lens is an approach to establish design specifications which generally define conditions of performance or dimensions that the optical system should meet.\nThe first step of this approach is to review previously developed optical systems that are similar to the optical system to be developed. From the conclusion of the review, a basic design is established and the target performance and the type of lens are determined. Then, the power arrangement is determined for the optical lens system by means of approximated calculation.\nWhile maintaining the determined power arrangement, the configuration is then determined as a thin-walled lens system in which the optical system is composed of lenses whose thicknesses can be neglected compared to the refractivity-dependent radius of curvature.\nThen, thicknesses are introduced in the configuration which has been determined as the thin-walled lens system and thus specific shapes of the lenses are determined. The design of a lens is carried out according to such a procedure described above.\nThe determination of the power arrangement is an approach carried out to satisfy the design specifications by regarding the total optical lens system as a combination of a plurality of optical subsystems and by determining the power of the individual optical subsystems and the distances between the principal points of the optical subsystems.\nIn the designing of a lens system, it is important to determine the specific shapes of the individual optical subsystems without breaking the power arrangement. To this end, an optical CAD system is used to evaluate and investigate the optical system, and to determine the specific shapes of the lenses.\nA zoom lens is an optical system in which the focal length of the total optical lens system can be continuously changed without changing the position of the image point, by changing the distance between the principal points of optical subsystems along the optical axis.\nThe distances between respective optical subsystems, which were determined when the power arrangement was determined, correspond to the distances between principal points of respective optical subsystems in the actual configuration.\nReferring to FIG. 41 and 42, the distance between optical subsystems 1 and 2, which was determined during the power arrangement, corresponds to the distance between the image-side principal point of optical subsystem 1 and the object-side principal point of optical subsystem 2 in the actual configuration.\nAs a result, if a zoom lens system is defined in an optical CAD system according to a conventional method of inputting the data of an optical lens system, a variable parameter is not given as a distance between optical subsystems as described above, but it is necessary to specify the numbers of the lens faces whose distances are to be varied and it is necessary to define a variable parameter as a distance between faces that belong to the face.\nTherefore, to determined a specific shape of a lens, the position of the principal point in this specific shape of the lens should be calculated and the parameter representing the distance between faces belonging to the face should by determined by specifying the number of the lens face to be varied. Thus, a very troublesome task is required to define the data of the optical system in an optical CAD system.\nFurthermore, because only one object position can be defined for one condition of an optical lens system, if an optical evaluation for different positions of the object is required for a certain focal length of the optical lens system, the data should be defined separately for each condition of the optical lens system corresponding to each object position. As a result, for each given object position, a condition should be calculated in which focusing is adjusted depending on the position of the object so that the conjugate plane with respect to the object position may coincide with the evaluating position of the image. Then, the number of the variable face to be varied for focusing should be specified and the distance between faces belonging to the face should be given as the result that was calculated in such a way as described above. Unless such a very troublesome task is carried out, the desired evaluation cannot be accomplished.\nFor example, to define the respective conditions shown in FIGS. 43 and 44, the data representing each condition of the optical system should be separately given to perform an optical evaluation.\nFIG. 43 shows a condition of an optical lens system where focusing is adjusted to an infinite object position. FIG. 44 shows a condition of the optical lens system where focusing is adjusted to an object position which is 20 cm distant.\nIn the case of an optical lens system of a camera shown in FIGS. 45 and 46, to make an optical evaluation of this optical system with respect to the two different positions of an object to be photographed, i.e., infinite and short distance positions, a focusing condition should be calculated in which the conjugate plane of the object position may coincide with the evaluating position of the image, i.e., the image is in focus, and then the data should be input separately for each of two conditions so that the data may represent the conditions shown in FIGS. 45 and 46.\nIn the case of the optical system shown in FIGS. 45 and 46, focusing is carried out by moving the position of the lens that is first when counted from the object side.\nIn some cases, an optical system comprises a large number of lens faces such as a few tens of lens faces, and it would be a very troublesome task for a user to input position data for each face, and input errors would often occur. Even if only a portion of the data representing the 3-dimensional positions of the respective faces of a group of optical elements is different, the data similar to each group would have to be input repeatedly.\nTo determine the spatial location for each lens face of a decentered optical system for each of the different zooming conditions, a memory area should be designated for the spatial locations corresponding to the number of different zooming conditions for each of the lens faces, which requires an unnecessarily large memory capacity. Thus, if a large number of zooming conditions are required, then the memory area for spatial locations should be expanded."}
{"text": "Mobile communication networks deploy various techniques for measuring the geographical locations of communication terminals. Such techniques are used, for example, for providing Location Based Services (LBS) and emergency services in cellular networks.\nSome communication interception systems masquerade as base stations and log the International Mobile Subscriber Identity (IMSI) numbers of communication terminals in their area. Such systems are commonly known as “IMSI catchers” or “active off-air interrogation systems.”\nPCT International Publication WO 2010/116292, whose disclosure is incorporated herein by reference, describes a system for obtaining information relating to an idle mobile station in a cellular network. The system generates and sends at least one input signal through a radio network controller to a radio access network, and identifies in data outputted by the radio network controller an output signal resulting from the input signal.\nNetwork Measurement Report (NMR) is a process in which a mobile communication terminal measures and stores signal strengths of neighboring base stations. NMR information is used, for example, for determining when to hand-off a mobile station from one base station to another."}
{"text": "The present invention pertains to a process which can be fully automated for accurately determining the alleles of STR genetic markers. More specifically, the present invention is related to performing PCR amplification on DNA, assaying the PCR products, and then determining the genotype of the PCR products. The invention also pertains to systems which can effectively use this genotyping information.\nTo study polymorphisms in genomes, reliable allele determination of genetic markers is required for accurate genotyping. A genetic marker corresponds to a relatively unique location on a genome, with normal mammalian individuals having two (possibly identical) alleles 104 for a marker on an autosomal chromosome 102, referring to FIG. 1A. (Though there are other cases of 0, 1, or many alleles that this invention. addresses, this if characterization suffices for the background introduction.). One important class of markers is the CA-repeat loci. This class is abundantly represented throughout the genomes of many species, including humans.\nA CA-repeat marker allele is comprised of a nucleic acid word 106\nPQRST,\nwhere P is the left PCR primer, T defines the right PCR primer, Q and S are relatively fixed sequences, and the primary variation occurs in the sequence R, which is a tandemly repeated sequence 108 of the dinucleotide CA, i.e.,\nR=(CA)n,\nwhere is n is an integer that generally ranges between ten and fifty. Thus, the length of the allele sequence uniquely determines the content of the sequence, since the only polymorphism is in the length of R.\nOne can therefore obtain genomic DNA, perform PCR amplification of a CA-repeat genetic marker location, and then assay the length of the allele sequences by differential sizing, typically done by differential migration of DNA molecules using gel electrophoresis. The resulting gel 110 should, in principle, fib clearly show the alleles of marker for each individual\"\"s genome. Further, these sizes can be determined quantitatively by reference to molecular weight markers 112.\nHowever, the PCR amplification of a CA-repeat location produces an artifact, often termed xe2x80x9cPCR stutterxe2x80x9d. Most likely due to slippage of the polymerase molecule on the nucleic acid polymer go in the highly repetitive CA-repeat region, the result is that PCR products are produced that correspond to deletions of tandem CA molecules in the repeat region. Thus, instead of a single band on a gel corresponding to the one molecule\nPQ(CA)nST,\nan entire population of different size bands\n{PQ(CA)nST, PQ(CA)nxe2x88x921ST, PQ(CA)nxe2x88x922ST, . . . }\nin varying concentrations is observed. This PCR stuttering 114 can be viewed as a spatial pattern p(x), or, alternatively, as a response function r(t) of an impulse signal corresponding to the assayed allele.\nThe stutter artifact can be extremely problematic when the two alleles of an autosomal CA-repeat marker are close in size. Then, their two stutter patterns overlap, producing a complex signal 116. In the presence of background measurement noise, this complexity often precludes unambiguous determination of the two alleles. To date, this has prevented reliable automated (or even manual) genotyping of CA-repeat markers from differential sizing assays.\nThis overlap of stutter patterns can be modeled as a superposition of two corrupted signals. Importantly, (1) the corrupting response function is roughly identical for two closely sized alleles of the same CA-repeat marker, and (2) this response function is largely determined by the specific CA-repeat marker, the PCR conditions, and possibly the relative size of the allele. Thus, the response functions 114 can be assayed separately from the genotyping experiment 116. By combining 118 the corrupted signal together with the determined response functions of the CA-repeat marker, the true uncorrupted allele sizes can be determined, and reliable genotyping can be performed.\nA primary goal of the NIH/DOE Human Genome Project during its initial 5 year phase of operation was to develop a genetic map of humans with markers spaced 2 to 5 cM apart (E. P. Hoffman, xe2x80x9cThe Human Genome Project: Current and future impact,xe2x80x9d Am. J. Hum. Genet., vol. 54, pp. 129-136, 1994), incorporated by reference. This task has already been largely accomplished in half the time anticipated, with markers that are far more informative than originally hoped for. In these new genetic maps, restriction fragment length polymorphism (RFLP) loci have been entirely replaced by CA repeat loci (dinucleotide repeats, also termed xe2x80x9cmicrosatellitesxe2x80x9d) (J. Weber and P. May, xe2x80x9cAbundant class of human DNA polymorphisms which can be typed using the polymerase chain reaction,xe2x80x9d Am J Hum Genet, vol. 44, pp. 388-396, 1989; J. Weber, xe2x80x9cLength Polymorphisms in dC-dA . . . dG-dT Sequences,xe2x80x9d Marshfield Clinic, Marshfield, Wis., assignee code 354770, U.S. Pat. #5,075,217, 1991), incorporated by reference, and other short tandem repeat markers (STRs). It is expected that at least 30,000 CA-repeat markers will be made available in public databases in the form of PCR primer sequences and reaction conditions. One of the advantages of CA repeat loci is their high density in the genome, with about 1 informative CA repeat every 50,000 bp: this permits a theoretical density of approximately 20 loci per centimorgan. Another advantage of CA repeat polymorphisms is their informativeness, with most loci in common use having PIC values of over 0.70 (J. Weissenbach, c. Gyapay, C. Dib, A. Vignal, J. Morissette, P. Millasseau, G. Vaysseix, and M. Lathrop, xe2x80x9cA second generation linkage map of the human genome,xe2x80x9d Nature, vol. 359, pp. 794-801, 1992; G. Gyapay, et. al., Nature Genetics, vol. 7, pp. 246-239, 1994), incorporated by reference. Finally, these markers 43; are PCR-based, permitting rapid genotyping using minute quantities of input genomic DNA. Taken together, these advantages have facilitated linkage studies by orders of magnitude: a single full-time scientist can cover the entire genome at a 10 cM resolution and map a disease gene in an autosomal dominant disease family in about 1 year (D. A. Stephan, N. R. M. Buist, A. B. Chittenden, K. Ricker, J. Zhou, and E. P. Hoffman, xe2x80x9cA rippling muscle disease gene is localized to 1q41: evidence for multiple genes,xe2x80x9d Neurology, in press, 1994), incorporated by reference.\nThe CA repeat-based genetic maps are not without disadvantages. First, alleles are detected by size differences in PCR products, which often differ by as little as 2 bp in a 300 bp PCR product. Thus, these alleles must be distinguished using high resolution sequencing gels, which are more labor intensive and technically demanding to use than most other electrophoresis systems. Second, referring to FIG. 2, CA repeat loci often show secondary xe2x80x9cstutterxe2x80x9d or xe2x80x9cshadowxe2x80x9d bands in addition to the band corresponding to the primary allele, thereby complicating allele interpretation. These stutter bands may be due to errors in Taq polymerase replication during PCR, secondary structure in PCR products, or somatic mosaicism for allele size in a patient. Allele interpretation is further complicated by the differential mobility of the two complementary DNA strands of the PCR products when both are labelled. Finally, sequencing gels often show inconsistencies in mobility of DNA fragments, making it difficult to compare alleles of individuals between gels and often within a single gel. The most common experimental approach used for typing CA repeat alleles involves incorporation of radioactive nucleotide precursors into both strands of the PCR product. The combined consequence of stutter peaks and visualization of both strands of alleles differing by 2 bp often leads to considerable xe2x80x9cnoisexe2x80x9d on the resulting autoradiograph xe2x80x9csignalsxe2x80x9d, referring to FIG. 2, which then requires careful subjective interpretation by an experienced scientist in order to determine the true underlying two alleles.\nThe stuttered signals of di-, tri-, tetra-, and other polynucleotide repeats can be modeled as the convolution of the true allele sizes with a stutter pattern p(x). Under this model, the complex quantitative banding signal q(x) observed on a gel can be understood as the summation of shifted patterns p(x), with one shifted pattern for each allele size. A key fact is that generally only one p(x) function is associated with a given genetic marker, its PCR primers and conditions, and the allele size. In the important case of two alleles, where the two allele sizes are denoted by s and t, one can write the expression\nq(x)=(x3+x1), p(x).\nThe multiplication of the polynomial expressions (x3+x1) and p(x) is one implementation of the underlying (shift and add) convolution process. Given the observed data q(x) and the known stutter pattern p(x), one can therefore determine the unknown allele sizes s and t via a deconvolution procedure. (Note that this convolution/deconvolution model extends to analyses with more than two alleles.)\nA corollary of highly dense and informative genetic maps is the need to accurately acquire, analyze and store large volumes of data on each individual or family studied. For example, a genome-wide linkage analysis on a 30 member pedigree at 10 cM resolution would generate data for approximately 30,000 alleles, with many markers showing five or more alleles. Currently, alleles are visually interpreted and then manually entered into spreadsheets for analysis and storage. This approach requires a large amount of time and effort, and introduces the high likelihood of human error. Moreover, future studies of complex multifactorial disease loci will require large-scale genotyping on hundreds or thousands of individuals. Finally, manual genotyping is arduous, boring, time consuming, and highly error prone. Each of these features suggests that automation of genotype data generation, acquisition, interpretation, and storage is required to fully utilize the developing genetic maps. Some effort has been made to assist in allele identification and data storage (ABI Genotyper manual and software, Applied Biosystems Inc.), incorporated by reference. However, this software still requires substantial user interaction to, place manually assigned alleles into a spreadsheet, and is unable to deconvolve (hence cannot accurately genotype) closely spaced alleles or perform other needed analyses. Importantly, no essential use is made of a CA-repeat marker\"\"s PCR stutter response pattern by the ABI software or by any other disclosed method or system for genotyping.\nThe Duchenne/Becker muscular dystrophy (DMD/BMD) gene locus (dystrophin gene) (A. P. Monaco, R. L. Neve, C. Colletti-Feener, C. J. Bertelson, D. M. Kurnit, and L. M. Kunkel, xe2x80x9cIsolation of candidate cDNAs for portions of the Duchenne muscular dystrophy gene,xe2x80x9d Nature, vol. 323, pp. 646-650, 1986; M. Koenig, E. P. Hoffman, C. J. Bertelson, A. P. Monaco, C. Feener, and L. M. Kunkel, xe2x80x9cComplete cloning of the Duchenne muscular dystrophy cDNA and preliminary genomic organization of the DMD gene in normal and affected individuals,xe2x80x9d Cell, vol. 50, pp. 509-517, 1987), incorporated by reference, is a useful experimental system for illustrating the automation of genetic analysis. The dystrophin gene can be considered a mini-genome: it is by far the largest gene known to date (2.5 millionbase pairs); it has a high intragenic recombination rate (10 cM, i.e., 10% recombination between the 5xe2x80x2 and 3xe2x80x2 ends of the gene); and it has a considerable spontaneous mutation rate (10xe2x88x924 meioses). Mutation of the dystrophin gene results in one of the most common human lethal genetic diseases, and the lack of therapies for DMD demands that molecular diagnostics be optimized. The gene is very well characterized, with both precise genetic maps (C. Oudet, R. Heilig, and J. Mandel, xe2x80x9cAn informative polymorphism detectable by polymerase chain reaction at the 3xe2x80x2 end of the dystrophin gene,xe2x80x9d Hum Genet, vol. 84, pp. 283-285, 1990), incorporated by reference, and physical maps (M. Burmeister, A. Monaco, E. Gillard, G., van Ommen, N. Affara, M. Ferguson-Smith, L. Kunkel, and H. Lehrach, xe2x80x9cA 10-megabase physical map of human Xp21, including the Duchenne muscular dystrophy gene,xe2x80x9d Genomics, vol. 2, pp. 189-202, 1988), incorporated by reference. Finally, approximately one dozen CA repeat loci distributed throughout the dystrophin gene have been isolated and characterized (A. Beggs and L. Kunkel, xe2x80x9cA polymorphic CACA repeat in the 3xe2x80x2 untranslated region of dystrophin,xe2x80x9d Nucleic Acids Res, vol. 18, pp. 1931, 1990; C. Oudet, R. Heilig, and J. Mandel, xe2x80x9cAn informative polymorphism detectable by polymerase chain reaction at the 3xe2x80x2 end of the dystrophin gene,xe2x80x9d Hum Genet, vol. 84, pp. 283-285, 1990; P. Clemens, R. Fenwick, J. Chamberlain, R. Gibbs, M. de Andrade, R. Chakraborty, and C. Caskey, xe2x80x9cLinkage analysis for Duchenne and Becker muscular dystrophies using dinucleotide repeat polymorphisms,xe2x80x9d Am J Hum Genet, vol. 49, pp. 951-960, 1991; C. Feener, F. Boyce, and L. Kunkel, xe2x80x9cRapid detection of CA polymorphisms in cloned DNA: application to the 5xe2x80x2 region of the dystrophin gene,xe2x80x9d Am J Hum Genet, vol. 48, pp. 621-627, 1991), incorporated by reference.\nMany of the problems with interpretation of dystrophin gene CA repeat allele data can be overcome by single or multiplex fluorescent PCR and data acquisition on automated sequencers (L. S. Schwartz, J. Tarleton, B. Popovich, W. K. Seltzer, and E. P. Hoffman, xe2x80x9cFluorescent Multiplex Linkage Analysis and Carrier Detection for Duchenne/Becker Muscular Dystrophy,xe2x80x9d Am. J. Hum. Genet., vol. 51, pp. 721-729, 1992), incorporated by reference. This approach uses fluorescently labeled PCR primers to simultaneously amplify four CA repeat loci in a single reaction. By visualizing only a single strand of the PCR product, and by reducing the cycle number, much of the noise associated with these CA repeat loci was eliminated. Moreover, the production of fluorescent multiplex reaction kits provides a standard source of reagents which do not deteriorate for several years following the fluorescent labeling reactions. In this previous report, referring to FIG. 2, alleles were manually interpreted from the automated sequencer traces. Coverage of the entire human genome at 10 cM resolution in fluorescently labeled polynucleotide markers for use in semiautomated genotyping is available. (Map Pairs, Research Genetics, Huntsville, Ala.; P. W. Reed, J. L. Davies, J. B. Copeman, S. T. Bennett, S. M. Palmer, L. E. Pritchard, S. C. L. Gough, Y. Kawaguchi, H. J. Cordell, K. M. Balfour, S. C. Jenkins, E. E. Powell, A. Vignal, and J. A. Todd, xe2x80x9cChromosome-specific microsatellite sets for fluorescence-based, semi-automated genome mapping,xe2x80x9d Nature Genetics, in press, 1994), incorporated by reference.\nThis invention pertains to automating data acquisition and interpretation for any STR genetic marker. In the preferred embodiment, the invention: identifies each of the marker alleles at an STR locus in an organism; deconvolves complex xe2x80x9cstutteredxe2x80x9d alleles which differ by as few as two bp (i.e., at the limits of signal/noise); makes this genotyping information available for further genetic analysis. For example, to establish DMD diagnosis by linkage analysis in pedigrees, the application system: identifies each of the dystrophin gene alleles in pedigree members; deconvolves complex xe2x80x9cstutteredxe2x80x9d alleles which differ by only two bp where signal/noise is a particular problem; reconstructs the pedigrees from lane assignment information;, sets phase in females; propagates haplotypes through the pedigree; identifies female carriers and affected males in the pedigree based on computer derivation of an at-risk haplotype; detects and localizes recombination events within the pedigree. Other uses of automatically acquired STR genetic marker data are the construction of genetic maps (T. C. Matise, M. W. Perlin, and A. Chakravarti, xe2x80x9cAutomated construction of genetic linkage maps using an expert system (MultiMap): application to 1268 human microsatellite markers,xe2x80x9d Nature Genetics, vol. 6, no. 4, pp. 384-390, 1994), incorporated by reference, the localization of genetic traits onto chromosomes (J. Ott, Analysis of Human Genetic Linkage, Revised Edition. Baltimore, Maryland: The Johns Hopkins University Press, 1991), incorporated by reference, and the positional cloning of genes derived from such localizations (B.-S. Kerem, J. M. Rommens, J. A. Buchanan, D. Markiewicz, T. K. Cox, A. Chakravarti, M. Buchwald, and L.-C. Tsui, xe2x80x9cIdentification of the cystic fibrosis gene: genetic analysis,xe2x80x9d Science, vol. 245, pp. 1073-1080, 1989; J. R. Riordan, J. M. Rommens, B.-S. Kerem, N. Alon, R. Rozmahel, Z. Grzelczak, J. Zielenski, S. Lok, N. Plavsic, J.-L. Chou, M. L. Drumm, M. C. Iannuzzi, F. S. Collins, and L.-C. Tsui, xe2x80x9cIdentification of the cystic fibrosis gene: cloning and characterization of complementary DNA,xe2x80x9d Science, vol. 245, pp. 1066-1073, 1989), incorporated by reference.\nThe present invention pertains to a method for genotyping. The method comprises the steps of obtaining nucleic acid material from a genome. Then there is the step of amplifying location of the material. Next there is the step of assaying the amplified material based on size and concentration. Then there is the step of converting the assayed amplified material into a first set of electrical signals corresponding to size and concentration of the amplified material at the location. Then there is the step of operating on the first set of electrical signals produced from the amplified material with a second set of electrical signals corresponding to a response pattern of the location to produce a third set of clean electrical signals corresponding to the size and multiplicities of the unamplified material on the genome at the location.\nThe present invention also pertains to a system for genotyping. The system comprises means or a mechanism for obtaining nucleic acid material from a genome. The system also comprises means or a mechanism for amplifying a location of the material. The amplified means or mechanism is in communication with the nucleic acid material. Additionally, the system comprises means or a mechanism for assaying the amplified material based on the size and concentration. The assaying means or mechanism is in communication with the amplifying means or mechanism. The system moreover comprises means or a mechanism for converting the assayed amplified material into a first set of electrical signals corresponding to size and concentration of the amplified material at the location. The converting means or mechanism is in communication with the assaying means. The system for genotyping comprises means or a mechanism for operating on the first set of electrical signals produced from the amplified material with a second set of electrical signals corresponding to a response pattern of the location to produce a third set of clean electrical signals corresponding to the size and multiplicities of the unamplified material on the genome at the location. The operating means or mechanism is in communication with the sets of electrical signals. The present invention also pertains to a method of analyzing genetic material of an organism. The present invention additionally pertains to a method for producing a gene."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of reed switch testing and in particular to an improved method and apparatus for use in detecting non-parallel reed switch contacts by measuring the duration or pulse width of each contact bounce after the initial operation or closure of the reed switch contacts and comparing each open contact duration to a programmed reference number corresponding to an acceptable duration.\n2. Description of the Prior Art\nIn the manufacturing of reed switch devices it is necessary to identify reed switches that have non-parallel contacts, generally called bad mating parts. These non-parallel contact switches usually have one reed member or blade twisted axially which produces a line contact instead of the desired surface contact between the reed members. This undesirable line contact gives a high but acceptable electrical contact resistance and yet may still cause an ultimate failure in the operation of the switch after operating in the field.\nVarious testing methods have been developed to detect these potential failure parts including the visual sorting one-by-one which of course includes the inherent disadvantage of the varying visual perception of the testing personnel. Electrical testing apparatus has been developed that applies a linearly changing current simultaneously to an operating coil and to its reed blades and then differentiating the voltage developed across the contacts to determine the operating characteristics of the switch such as those of the type described in U.S. Pat. No. 3,586,961, which issued to R. M. Rovnyak on June 22, 1971 and U.S. Pat. No. 3,822,398, which issued to R. M. Rovnyak on July 2, 1974. These techniques however do not make use of the contact bounce time to determine badly mated contacts.\nOther relay testing circuits of the prior art measure various parameters such as the length of time that the contacts continue to chatter after initial closure, the number of bounce pulses present after closure, and the accumulative bounce time such as disclosed in U.S. Pat. No. 2,236,157, which issued to W. T. Rea on Mar. 25, 1941, U.S. Pat. No. 1,677,157, which issued to A. Weaver on July 17, 1928, U.S. Pat. No. 3,612,985, which issued to A. F. Rockett on Oct. 12, 1971, and U.S. Pat. No. 3,241,063, which issued to C. L. Beattie et al. on Mar. 15, 1966.\nOther prior art circuits have been devised to measure the pulse duration of a contact bounce as compared to a predetermined level by the use of capacitor charge and discharges such as disclosed in U.S. Pat. No. 3,364,419, which issued to G. D. Anderson on Jan. 16, 1968 and in an article entitled \"Relay Contact Bounce Measurements\", by Richard Jiu in Electronics Magazine, August, 1955 issue. Such circuits however do not include digital techniques that allow for the programming of various reference time durations. Further, this type of RC time constant measurement is also less accurate than those obtained by digital techniques."}
{"text": "The Global Positioning System or “GPS” is by definition ubiquitous and is now widely used for a wide variety of military, commercial and consumer applications. GPS is used for everything from missile guidance to guiding a backpacker who has wandered off the Pacific Crest Trail in a surprise Sierra Nevada snowstorm. Details of GPS system operation are published elsewhere. Suffice to say here that the GPS network of satellites provides information to a GPS receiver as long as it is “in view” of a least a handful of the dozen or so GPS satellites positioned in geostationary orbit. The GPS receiver or an associated processor calculates the current location of the receiver from the received information. The resulting location data may be expressed for example as latitude, longitude, altitude or in an equivalent coordinate system.\nIn recent years, the size and cost of GPS receivers and related equipment have diminished so dramatically so that the technology is now available in packages small enough to hand carry and at pricing levels that bring the equipment into the consumer marketplace. In fact, many mobile phones today include a built-in GPS receiver.\nOne growing application of GPS receivers is for navigation purposes in motor vehicles. Where a driver is unfamiliar with a desired destination location or the surrounding area, GPS-based navigation equipment can help, as it can determine the location of the vehicle at almost any time. Commercial vehicle applications are growing as well, from local garbage trucks to long-haul rigs.\nOne vendor of GPS equipment, known as Garmin, recited in a press release: “Thanks to the real-time information transmitted through the GPS-Buddy unit, fleet managers can know where their fleet is and communicate with the drivers 24-hour-a-day, 7-days-a-week . . . . Every 60-seconds, a secured GPS-Buddy server receives information about the vehicle's location, speed, and status—such as engine data, loading door, and thermostat. The information is then available on the internet to GPS-Buddy customers. The data may be viewed at anytime on any device with internet access such as personal computers, Smartphones, or PDAs.” In short, it is known to upload location data and other information from a motor vehicle to a server.\nWhile vehicle (or hand carried) navigation systems generally can acquire their present location via GPS, they must in all cases be programmed with a desired destination location (hereafter simply, “destination”) before the navigation system can assist in navigating to that location.\nProgramming a navigation system with a desired destination may be difficult for some users who are unfamiliar with the user interface at hand, and at any rate takes the user's time.\nSystems and methods also are known for transmitting digital data over the voice channel of a wireless telecom network. Voice services have the advantages of low cost, high reliability and wide availability across various wireless networks and technologies. These are sometimes called “concierge” services in view of the human contact (via phone) and personalized service. Digital wireless data services by contrast are sometimes unreliable, and can vary in bandwidth, delay and other parameters across different networks and technologies.\nMoreover, transmitting data in the voice channel has the characteristic that a voice call connection must be established. This enables substantially simultaneous voice and data communications. Thus, for example, an emergency call taker or concierge operator can talk to a person who requires assistance and, at substantially the same, time receive data such as the person's location or physiological or medical data. Live concierge operators can give verbal “turn-by-turn” driving instructions to guide a user to a desired location while the user is en route, typically by motor vehicle, as long as the vehicle is suitable equipped with integral wireless phone, voice-band or “in-band” signaling modem, etc. . . . Location data can be provided through the voice channel from the car to the concierge or other “call taker.” These services are expensive, however, because they are labor intensive. Some users prefer a more independent approach, using on-board electronic navigation equipment.\nOn-board navigation equipment is now available in some markets and built into some luxury automobiles and rental cars. These on-board navigation devices generally are “stand-alone” meaning that they operate without the assistance of a remote, live operator. They typically have a GPS receiver built into the equipment, or connected to it, to determine a present location of the car. They also have databases (e.g. stored on CD-ROM) with maps or other GIS data to aide navigation, or they can download them as required, for example via the wireless telecommunications network, internet-connected WAP, satellite etc. The stand-alone systems may provide recorded or simulated voice announcement of so-called “turn-by-turn” driving instructions to guide the user to a selected destination. In all cases the desired destination location must be input to the system. Programming a navigation system, i.e., to input the desired destination location, is clumsy and difficult, for example because of a very limited user interface, small display screen, etc. It is especially difficult and may be dangerous to attempt to input a desired destination address to a navigation device while a user is driving."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a zoom lens and, more particularly, to a compact wide-angle zoom lens suitable for use with a 35 mm format camera.\n2. Related Background Art\nIn recent years, a super wide-angle lens and a wide-angle lens employed as an interchangeable lens for a 35 mm format still camera have become smaller and smaller in size and higher and higher in performance. Especially for actualizing a small-sized wide-angle zoom lens at a low price, the optimum lens type is a so-called negative/positive two-group zoom lens. A variety of types have hitherto been proposed. Further, it has recently been possible to attain a miniaturization, a higher performance and a drop in costs by incorporating an aspherical lens surface. Examples of such known zoom lenses are shown in Japanese Patent Application Laid-open Nos. 62-94812 and 63-136013, and U.S. Pat. No. 5,076,677.\nThe wide-angle zoom lens disclosed in Japanese Patent Application Laid-open No. 62-94812 stated above, however, presents the following defect. This wide-angle zoom lens is, though small in the number of constructional lenses and also relatively small in total thickness of the lens system, relatively large in total length at the wide-angle end. A diameter of the front lens element increases. It is difficult to install a .phi.52 mm filter that has already been generalized as a standard filter size. The following is an additional defect. If forcibly installed, this brings about a remarkable reduction in amount of marginal rays at the wide-angle end. The principal ray having the maximum field angle is eclipsed during near distance photographing. Besides, in an aspect of aberrations also, it is insufficient to compensate a coma, an astigmatism and a curvature of field in the vicinity of the maximum field angle at the wide-angle end. It is also insufficient to compensate a chromatic aberration of magnification and a downward coma (hereafter called a [downward chromatic coma]) due to a difference between wavelengths of rays.\nThe super wide-angle zoom lens disclosed in Japanese Patent Application Laid-open No. 63-136013 stated above exhibits such an advantage that this zoom lens has the maximum field angle as large as 92.degree. or more. The super wide-angle zoom lens is, however, constructed of a Great number of lenses and is large both in size and in diameter of the front lens element. This zoom lens is therefore insufficient in terms of reducing the size and the costs as well. A filter size is also large. Further, in the aspect of aberrations, the compensations of the curvature of field and the astigmatism are not enough from the wide-angle end to a state of intermediate focal length.\nIn addition, U.S. Pat. No. 5,076,677 by the same applicant as the one of the present invention discloses a small-sized wide-angle zoom lens having a field angle of approximately 76.degree. at the wide angle and a small diameter of the front lens element. If the field angle at the wide-angle end increases up to as large as 84.degree. with this lens construction, a total length at the wide-angle end outstandingly elongates. The diameter of the front lens element is also dilated, which in turn causes an increase in the filter size. Further, the amount of marginal rays is conspicuously reduced. Eventually, this conduces to a defect in which the eclipse occurs. Besides, in terms of aberrations, there are undesirably worsened particularly the curvature of field, the astigmatism, the downward coma, the downward chromatic coma and the distortion.\nHence, it is difficult to actualize the small-sized wide-angle zoom lens exhibiting a good performance with such refracting power arrangements and lens constructions of the known zoom lenses."}
{"text": "The present invention relates to pharmaceutical tablet compositions that contain crystalline aripiprazole Type II and pharmaceutically acceptable excipients.\nAripiprazole or more properly 7-[4-[4-(2,3-dichlorophenyl)-1-piperazinyl]butoxy]-3,4-dihydrocarbostyril, is a compound of the formula (1).\nIt is a commercially marketed, pharmaceutically active substance useful for treatment of schizophrenia and bipolar disorder. It is disclosed in EP 367141/U.S. Pat. No. 5,006,528. The commercially marketed product is the free base of the title compound (1).\nSolid state aripiprazole was prepared in U.S. Pat. No. 5,006,528 by a two-fold recrystallization of crude aripiprazole from ethanol resulting in colorless flake crystals having a melting point of 139-139.5° C. In an article of Aoki (Study on Crystal Transformation of Aripiprazole, The Fourth Japan-Korea Symposium on Separation Technology, p.937 ff (1996)), this solid state form was designated as Type I aripiprazole and identified as an anhydrate. Aoki also teaches that the Type I aripiprazole may be converted into a Type II aripiprazole by heating at 130-140° C. for 15 hours. This product is also an anhydrate with a melting point of 150° C. When both Type I and Type II aripiprazole were recrystallized from an alcoholic solvent containing water up to 20%, the product was an aripiprazole hydrate labeled as Type III by Aoki. Type III aripiprazole can be converted into Type I by heating at 80° C.\nWO 03/26659 (EP 1330249) teaches that Type I aripiprazole, the alleged original solid form of aripiprazole, is significantly hygroscopic. In an effort to find a form of aripiprazole having reduced hygroscopicity and better processing qualities, seven crystalline forms (A-G) were described.\nHydrate Form A is taught as a useful intermediate for making anhydrate forms. Hydrate Form A can be prepared by milling Aoki's hydrated Type III. Contrary to the conventional Form III hydrate, the Hydrate Form A does not exhibit sharp dehydration endothermic peak at 123.5° C. at TGA, but has a gradual endothermic peak between 60-120° C.\nAnhydrous Form B, which seems to be the preferred crystalline form, is not hygroscopic; i.e., less than 0.4% water uptake in 24 hours, and is a stable crystalline form. It can be prepared by heating the Hydrated Form A, preferably at 90-125° C. for 3-50 hours or by heating the Type I/Type II aripiprazole at 90-125° C. Although the Anhydrate Form B of WO 03/26659 is not hygroscopic, it suffers from being unsuitable for milling. Specifically, if milling is attempted in order to create small particle sizes such as 50 microns or less, the milled substance tends to adhere to the milling machine making an industrial process difficult. To overcome this problem, WO 03/26659 teaches forming Hydrate Form A aripiprazole, milling the Hydrate Form A to the desired size and then heat converting to Anhydrate Form B.\nThe other anhydrate forms disclosed therein are briefly summarized below:\nForm C: Prepared by heating an aripiprazole anhydrate to 140-150° C. Endothermic peak around 150.2° C.\nForm D: Prepared by recrystallization of aripiprazole anhydrate from toluene. Endothermic peaks at 136.8 and 141.6° C.\nForm E: Prepared by double heating, dissolving, and crystallizing aripiprazole in acetonitrile with crystallization at about 70° C. Endothermic peak at 146.5° C.\nForm F: Prepared by heating a suspension of aripiprazole anhydrate in acetone. Endothermic peaks at 137.5 and 149.8° C.\nForm G: Prepared by putting glassy state of aripiprazole anhydrate in a sealed vessel and keeping it at room temperature for at least 2 weeks. Exothermic peak at 122.7° C., endothermic peak at 141.0° C.\nMore recently techniques for crystallizing Type II aripiprazole directly from a solution were disclosed in U.S. Provisional application Serial Nos. 60/662,552, filed Mar. 17, 2005, and 60/692,557, filed Jun. 22, 2005; the entire contents of each being incorporated herein by reference. In general the Type II was found to be obtainable by precipitation from 2-propanol, dimethyl sulfoxide, or tetrahydrofuran, with a mixture of 2-propanol and ethyl acetate being a preferred embodiment, especially at somewhat elevated temperatures. Similarly, WO2005/058835 also purports to have directly crystallized Type II aripiprazole.\nAlthough various polymorphs of crystalline aripiprazole are known, it is believed that the commercial aripiprazole product uses Type I aripiprazole. However, the commercial product apparently suffers from moisture sensitivity in that the U.S. product is packaged with a large amount of desiccant (3 mg instead of the usual 2 mg desiccant) and the EU product is packaged in moisture protecting aluminum-aluminum blisters.\nIt would be desirable to form an aripiprazole tablet composition that is less moisture sensitive than the current commercial product."}
{"text": "The invention relates to fabrication of thin film oxide superconductors.\nMetal organic deposition (MOD) is a method for fabrication of uniform thickness thin films from liquid precursor solutions as described by Y. L. Chen et al., J. Mater. Res., 4, 1065 (1989). Commonly, metal carboxylates dissolved in organic solvents are deposited on polished substrates by dipping or spin coating with a photoresist spinner. The green precursor films produced by this coating process are transformed into superconductor films by selective heat treatment which often involves multiple distinct calcination steps. Metal organic deposition techniques can be used to deposit films of uniform and controlled composition rapidly and without a need for sophisticated vacuum systems, J. V. Mantese et al., MRS Bull., 14, 48 (1989).\nSeveral investigators have reported preparation of Ba.sub.2 YCu.sub.3 O.sub.7-x (BYC) thin films using metal organic deposition including, Y-M. Chiang et al., Mat. Res. Soc. Symp. Proc., 99, 307 (1988); M. E. Gross et al., Appl. Phys. Lett., 52, 160 (1988); T. Kumagai et al., Chem. Lett., 1645 (1987); M. Klee et al., J. Crys. Growth, 91, 346, (1988). However, these workers observe that the electrical transport properties of metal organic deposited BYC films prepared from metal carboxylate precursors are generally inferior to sputtered, coevaporated and laser ablated films such as those prepared by R. L. Burton et al., submitted to SERI conference, Colorado Springs, Colo., November, 1988; S-W. Chan et al., J. Appl. Phys., 65, 4720 (1989); and P. M. Mankiewich et al., Appl. Phys. Lett., 51, 1753 (1987). Residual precursor materials trapped in MOD films may be responsible for these inferior properties; however, the high processing temperatures needed to completely eliminate precursor residues can cause undesirable reactions between the superconductor film and substrate."}
{"text": "(a) Field of the Invention\nThe present invention relates to an electron gun assembly for a cathode ray tube (CRT) and, more particularly, to an electron gun well adapted for a monochrome CRT to be mounted within a projection display device to realize a monochrome image.\n(b) Description of the Related Art\nGenerally, a CRT-based projection display device mainly has three monochrome CRTs for realizing red (R), green (G) and blue (B) monochrome images, and an optical system for amplifying the monochrome images made at the three CRTs and projecting the amplified images to a projection screen to produce color images.\nAs the monochrome CRT scans the display screen with one stream of electron beams, and the respective monochrome CRT screen images are projected to the projection screen while being amplified by about ten times, the brightness of the display screen related to the monochrome CRT is lower than that of the display screen related to the usual CRT. Therefore, compared to the usual CRT, relatively high electric currents need to be applied to the electron gun for the monochrome CRT to heighten the brightness of the display screen.\nUsually, the electron gun for the monochrome CRT emits electron beams with the application of electric currents ranging from 0.5 mA to 3 mA, which are two or three times more than those applied to the electron gun for a usual color CRT ranging from 0.2 mA to 1 mA. A high unipotential focus (Hi-UPF) type exhibiting an excellent focus characteristic in the range of higher currents is commonly used for the monochrome CRT electron gun.\nWith the Hi-UPF type electron gun, the second electrode receives the screen voltage, and the fourth electrode receives the focus voltage. A third electrode is placed between the second electrode and the fourth electrode to receive a high anode voltage (roughly, 32 kV). A strong pre-focus lens is formed between the second and the third electrodes due to the high potential difference between the second and the third electrodes, and reduces the spot size of electron beams in the range of higher electric currents.\nFurther, the monochrome CRT electron gun serves to make formation of monochrome images practically under the application of electric currents of 2 mA or less. With the available electric current range of 0.5-3 mA, the Hi-UPF type electron gun exhibits an excellent focus characteristic in the higher current range of more than 2 mA. By contrast, in the relatively lower current range of 2 mA or less, it turns out that the spot size of electron beams is increased.\nThe increase in the beam spot size occurs because when the electron\nThe increase in the beam spot size occurs because when the electron beam current is lowered, the crossover point of the electron beams formed at the triode portion moves from the second electrode to the third electrode, and the emission power of the electron beams incident upon the pre-focus lens is weakened. Consequently, with the electric current range of 2 mA or less, the spot size of the electron beams is increased while deteriorating the resolution, resulting in unclear display images.\nIn order to reduce the beam spot size with the current range of 2 mA or less, it has been proposed that the size of the beam-guide hole formed at the first electrode be reduced. However, this reduction makes the area of electron emission for the cathode so small that the life span of the electron gun with the cathode is reduced.\nU.S. Pat. No. 4,271,374 discloses a CRT electron gun with a structure where the equivalent diameter of the main-focus lens (formed between the fourth electrode receiving the focus voltage and the fifth electrode receiving the anode voltage) is enlarged to increase the capacity thereof.\nHowever, with the above structure, as the CRT neck portion mounting the electron gun thereon is limited in its diameter, there is a limit in mechanically enlarging the opening diameter of the fourth and the fifth electrodes forming the main-focus lens. Such a limit is made because the electron gun formation electrodes need to be spaced apart from the inner surface of the neck portion by a predetermined distance to grant the withstand voltage characteristic to the CRT. Accordingly, as the opening diameter of the fourth and the fifth electrodes is established in a predetermined manner, it is difficult to achieve the desired electrode capacity in an effective manner."}
{"text": "There is frequently a need to provide RF sensors at specific points within the whitegoods appliances and industrial machinery to allow for more accurate control or for control at specific points within the appliance and the machinery.\nWires may be cumbersome for the end user thus detracting from the user experience. Wires may not be feasible given the temperatures, movement, or other hazardous conditions between the location of the sensing point and the power supply. While this can be avoided with battery powered sensors, this approach may result in a life time issue requiring the customer to change batteries that may be troublesome or even impossible."}
{"text": "Internally illuminated signs are customarily required by fire and building codes of various governmental agencies to identify an exit from a given area or to identify a direction to an exit. Such internally illuminated signs often include a housing with a stencil which has the letters forming the word \"EXIT\" cut into a stencil as a message. A directional indicating arrow is cut into the stencil at each end of the word \"EXIT\". The sign includes an illumination source internally of the housing to provide the requisite illumination for the message.\nEach sign is typically mounted on a structure. A given sign may be mounted from a ceiling or supporting structure wherein the top of the sign is connected to the supporting structure. Accordingly, wiring to the interior of the sign and to light source enters the housing through a top wiring port. The sign may be mounted on a vertical wall so that one end of the sign is attached to the vertical wall and wiring to the interior of the sign and to light source enters the housing through a wiring port in an end wall of the sign. The sign may also be mounted with an opposite end against a vertical wall so that the wiring enters the housing through a wiring port in the other end of the housing.\nIn order to make a sign adaptable for mounting as a top mounted or a right or left mounted sign, the sign manufacturer typically provides wiring ports on the top and two opposed ends of the sign. It is desirable to provide a plate closure for the wiring ports which are not in use to close off the interior of the housing.\nA manufacturer typically provides removably closure plates for wiring apertures which are not in use. One of the problems encountered in the field by workmen is that the closure plates are difficult to handle. Therefore, it is particularly desirable to provide an improved removable closure plate for a wiring port, which closure plate may be handled with facility in the field."}
{"text": "The present invention relates to a plant for transmitting electric power between a direct voltage network for High Voltage Direct Current (HVDC), and at least two alternating voltage networks connected thereto through a station with each of the stations being adapted to perform transmission of electric power between the direct voltage network and a respective one of the alternating voltage networks.\nIn such plants, according to the prior art, primarily uninsulated overhead lines have been used to implement a direct voltage network interconnecting stations which may be located very large distances apart, for example, 1000 km. However, such overhead lines have a significant disturbing influence upon the country and the living things where they exist. The alternative has been to use cables having an inner conductor surrounded by a thick insulating layer formed by a paper impregnated by oil, but such cables are so costly that they are not a realistic alternative to overhead lines. Experiments have been carried out with cables having a conductor and an insulating layer on a polymer base surrounding the conductor for transmitting High Voltage Direct Current (HVDC), but these experiments have revealed great problems in using this type of cable for transmitting High Voltage Direct Current (HVDC), since space charges are created in the insulating layer which may cause breakthroughs and breakdowns in the cables. The space charges are caused by the changes of polarity which the conductor is subjected to when the feeding direction of the power transmitted between the direct voltage network and an alternating voltage network is changed. A result of these problems is that overhead lines have been used for transmitting High Voltage Direct Current (HVDC) over land, and cables having tailor made insulating sheets have been used as sea cables for transmitting between land areas separated by water."}
{"text": "1. Field of the Invention\nThis invention relates to interfacing a node, including an end node such as a work station, with a LAN (Local Area Network) link in a communication network. More particularly, this invention relates to maintaining data frame continuity without extensive data buffering in an environment where data rate capability of a link is much higher than the data rate capability of a node attached to the link.\n2. Description of Related Art\nOne technique for increasing the data capacity a work station may use in a LAN is to decrease the number of stations that share the data capacity of that LAN. In FIG. 1, work stations 10 and 12, requiring very high capacity (for example, file servers) can be attached to a dedicated switch port in a dedicated capacity concentrator 14. Concentrator 14 contains a data transfer unit 13, such as a crossbar switch and ports 15. Ports 15 connect to communication links that, in turn, are connected to work stations or other LAN networks. Ports 15 would contain extensive buffering and communication media controls to implement the communication media protocol.\nPorts 15A and 15B are connected as full duplex links to work stations 10 and 12, respectively. Port 15C is connected to shared capacity concentrator 16. Concentrator 16 has work stations 18, 19 and 20 connected to it in a ring network configuration. If the capacity of the ring network 16 is 100 Mbps (mega bits per second) and there are three work stations sharing the capacity of the ring, each station effectively has an average useable data rate of 33.3 Mbps. In contrast, each of the stations 10 and 12 has its own full duplex communication link, and uses, simultaneously, the entire 100 Mbps capacity in each direction of the full duplex link (effectively 200 Mbps capacity).\nA Local Area Network with single full duplex 100 Mbps links between crossbar switches in the network is taught in an article by M. D. Schroeder et al, entitled \"Autonet: A High-Speed, Self-Configuring Local Area Network Using Point-to-Point Links,\" published in IEEE Journal On Selected Areas In Communications, Vol. 9, No. 8, October 1991 (pp. 1318-1335).\nSchroeder et al, in the above-cited article, recognize that in a network with high capacity data links, the data may back up through the network if a crossbar switch cannot keep up with the data rate. Autonet provides 4 kB (kilo byte) FIFO (First In First Out) buffers at the receiving switch ports in the crossbar switches and flow control for the links. The flow control between sending and receiving crossbar switch stations is accomplished by controls in the receiving switch station issuing start/stop signals to the sending switch station.\nThis problem of data backup on high-speed data links is greatly accentuated at the link interface to user work stations. Usually, a work station has an internal bus that has a much lower capacity than the data link. For example, a personal computer, as a work station, has an internal bus with a typical capacity of approximately 20 Mbps in contrast to a high-speed data link with a capacity of 100-200 Mbps. In this situation, data frames must be buffered at the work station during both receive and transmit operations. Since data frames are lost if a buffer overflows or underflows, and since data frames are typically 4 kB (kilo Bytes) in length, large buffers, 64 kB to 256 kB, are required at the work station. In a communication network switching node, or in an end node (personal computer system or work station), this equates to having to add expensive memory chips on each communication adapter card."}
{"text": "1. Field\nThe present disclosure relates generally to high-speed data communications, and more particularly, to asymmetric communications between components of electronic devices and in particular to multi-phase encoding and decoding.\n2. Background\nHigh-speed interfaces are frequently used between circuits and components of mobile wireless devices and other complex apparatus. For example, certain devices may include processing, communications, storage and/or display devices that interact with one another through communications links. Some of these devices, including synchronous dynamic random access memory (SDRAM), may be capable of providing or consuming data and control information at processor clock rates. Other devices, such as display controllers, may require variable amounts of data at relatively low video refresh rates.\nHigh-speed interfaces are often limited by clock skew and are subject to interference. High frequency signals are often transmitted using differential interfaces to provide common-mode rejection for critical signals. In devices such as memory devices, which transmit and receive large amounts of data over wide data and control buses, interfaces maybe expensive and may consume significant power and real-estate on a circuit board."}
{"text": "Inventors have sought to develop artificial candles for many years that simulate actual candles. Many do this, however, replacing the lighting element of these artificial candles evades the skill of the average user. As a result, these expensive candles are thrown away, the prior art includes: U.S. Patent Application 2011/0280008 filed by Bakian; U.S. Pat. No. 7,744,266 issued to Higley; and U.S. Pat. No. 7,828,462 issued to Jensen.\nJensen and Bakian use a single light emitting diode (LED) attached to an integrated circuit board. Higley teaches a replaceable LED socket but does not explain how to mimic a wick. The present invention fills in the holes created by the prior art to teach a LED that is replaceable and can mimic a wick."}
{"text": "For a wide array of analytical techniques, accurate analysis of air or other gas contaminants requires the generation of standard atmospheres to calibrate and verify instrument response. Direct injection of a standard atmosphere into an analytical system can also be used as a diagnostic tool to evaluate the performance of individual system components and the overall system as a whole. This real-time diagnostic information can then be used to compensate for changes in instrument response due to changes in the operating environment or the deterioration of system components.\nPortable, handheld microanalytical systems, which have been termed “chemical laboratories on a chip,” are being developed to enable the rapid and sensitive detection of particular chemicals, including pollutants, high explosives, and chemical and biological warfare agents. These microanalytical systems should provide a high chemical selectivity, to discriminate against potential interferents, and the ability to perform the chemical analysis on a short time scale. These microanalytical systems also require accurate, reliable, and field-portable standards for calibration and evaluation of instrument response.\nA common analytical technique, particularly for microanalytical systems, is gas chromatography. A conventional gas chromatograph comprises a means for injection of a sample to be analyzed, a supply of carrier gas, a column for separating the components of the sample, and a detector for detecting the separated components that are eluted from the column. For quantitative analysis, it is necessary to accurately determine how much of a component is in the unknown sample. Generally, quantitative analysis requires the comparison of the size of the chromatographic peak of a component in an unknown sample to a chromatographic peak of the same component in a known calibration standard.\nIn addition to calibration, standards can be used to evaluate quantitative errors in the chromatographic process. Quantitative errors can be associated with peak size measurement, standardization, sampling techniques and sample introduction, and chromatographic system errors. In particular, standards can be used to identify column-induced changes in the sample character and detector errors due to overload or other factors. For example, a standard can be used to determine and correct for retention time shifts due to adsorption or other problems during gas chromatography measurements.\nBoth static and dynamic methods have been used to produce gas or vapor standards. Dynamic methods, based on continuous flow of a standard-containing gas mixture, have the advantages of reducing adsorption problems and being able to vary the concentration of the standard by simple change in the diluent flow. In general, the standard should be as close to the unknown sample as possible, stable, and of high purity. Especially with trace analysis, sample size of the standard and the unknown sample should be kept the same within measurement error.\nOne dynamic method to generate vapor standards is to use the diffusion of vapor through a capillary to add small amounts of the vapor to a flowing gas stream. In a conventional diffusion source, the liquid whose vapor is to provide the standard of interest is contained in a reservoir at a known temperature. The liquid is allowed to evaporate from the reservoir and the vapor diffuses through the capillary into a flowing diluent gas stream. The vapor concentration in the resulting gas mixture can then be determined from rate of diffusion of the vapor through the capillary and the flow-rate of the diluent gas.\nWhereas such diffusion tubes are well developed for conventional systems, a need exists for accurate, reliable, and field-portable standards for the calibration and evaluation of microanalytical systems. The present invention comprises a microfabricated diffusion source for use with such microanalytical systems. The invention can provide a very small sized calibration source that can be integrated with the fabrication of the microanalytical system by microelectromechanical systems (MEMS) technologies. This integrated fabrication eliminates the need for an external calibration source and also takes advantage of the economies of scale and low cost inherent with MEMS manufacturing. In addition, due to the very low calibrant flux requirements, which are typically in the nanograms or picograms per second, the lifetime of the microfabricated diffusion source can be months to years. These lifetimes can minimize, or in some cases eliminate, the need to periodically replenish the calibration source material."}
{"text": "1. Field of the Invention\nThe present invention relates to computer networks. More specifically, the present invention relates to a method and an apparatus for detecting a non-compliant router that improperly responds to ARP requests for IPv4 Link-Local addresses.\n2. Related Art\nIn order to participate in an Internet Protocol (IP) network, a network node, such as a computer or a printer, needs to be configured with an IP address. Additionally, in order communicate with another device using IP over Ethernet or similar media, the network node needs the ability to translate an IP address for the device into the corresponding hardware (e.g. Ethernet) address for the device. The Address Resolution Protocol (ARP) solves this translation problem by providing an address-resolution mechanism which determines the hardware address for a given IP address, thereby enabling a network node to participate in an IP network.\nIP addresses can be configured either manually by the user, or automatically with the help of another network node, such as a DHCP server. Unfortunately, a DHCP server may not always be available. Furthermore, it is cumbersome to have the user configure an IP address. Therefore, there is a strong need for a configuration mechanism by which a network node can automatically configure an IP address on its own.\nIPv4 Link-local addressing (RFC 3927) provides such a configuration mechanism. In IPv4 link-local addressing, a configuration mechanism within a network node randomly picks a “candidate” IP address in a specified link-local address range and sends an ARP request to check whether the candidate link-local IP address is unique within the scope of the link. If the candidate link-local IP address is already in use by another network node, then the configuration mechanism picks another link-local IP address and checks whether that is unique. Once the configuration mechanism finds a unique link-local IP address, the network node uses it to communicate with other nodes in the network.\nUnfortunately, some network routers are misconfigured to respond to all ARP requests, including ARP requests for any link-local IP address. If such a router is attached to a local network, it can interfere with the configuration mechanism by responding to ARP requests sent by the configuration mechanism, making it appear as if every single possible address is already in use. This can prevent the configuration mechanism from selecting an IP address.\nFor example, FIG. 1 illustrates how a router can interfere with the process of assigning an IP address. During operation, node A selects a candidate IP address (step 102). Next, node A probes for the candidate IP address by performing an ARP request (step 104). The interfering router then responds to the ARP request (step 106). This causes node A to infer that another node already has the IP address, and to select a different candidate IP address (step 108). The system then returns to step 104 and the process repeats indefinitely. Consequently, node A is unable to assign an IP address.\nNote that if such a router is configured to ignore ARP “probes,” which are sent by a device that is attempting to assign a link-local IP address, the router can still interfere with normal (non-probe) ARP requests, which are sent to determine an Ethernet address associated with an IP address.\nFor example, FIG. 2 illustrates how a router can interfere with normal ARP requests. First, a node A sends an ARP request with the IP address of a device that node A wants to communicate with (step 202). Next, a router incorrectly responds to the ARP request, interfering with the response from the actual device that node A wants to communicate with (step 204). This causes node A to obtain the wrong Ethernet address. IP packets are subsequently sent to the router instead of to the desired device, where they are lost because the router either forwards them incorrectly, or simply discards them without forwarding at all.\nHence, what is needed is a method and an apparatus for handling the above-described problems encountered while performing ARP requests."}
{"text": "There is prior art having to do with top-drive rigs with vertical pipe stacking U.S. Pat. No. 7,021,374 Weatherford ('374), as well as a system designed for snubbing and work-over rigs U.S. Pat. No. 6,158,516 Cudd Pressure Control ('516), and a more general pipe-racking system U.S. Pat. No. 4,042,123 Sheldon et al ('123).\nThere are, however, significant differences between these systems and this invention.\nFor instance, '374 Weatherford deals with a top-drive system, and thus has drive equipment which must align with the upper box-end of the tubing in the elevator, inside the drive unit, and the drive unit itself is integrated in the top-end. The similarity ends with the use of hydraulic rains connected to an off-set at the upper end of the bails and to the bails themselves lower down, to effect some vertical displacement. This is not exactly the same as this invention, and takes place in a completely different working environment, although the Weatherford's hydraulic offset device has some similarities, it is effected in a different manner, on a different rig type, with a different drive system.\nA block-retracting linkage between a vertical rail offset from the well's centre line and a travelling block is provided in '123 Sheldon, with a number (3) of hydraulic rams which can be actuated to move the block (and any suspended tubulars) off the well's centre line and toward and to a vertical pipe-rack means. This is different from the hydraulic bail/ram system of this invention in that the block in the system of this invention is suspended and not held by a rail-based mechanical setup.\nIn '516 Cudd, a combination coiled-tubing and rack-back jointed tubing rig is described, but without details of the elevator and racking system. This is relevant in that the type of tubular is commonly seen in pressure-controlled well settings, but Cudd does not disclose a hydraulic-ram bail swing system, and so is merely cited as a reference to similar rig environments in the prior art with rack-back tubular systems."}
{"text": "This invention relates to the attachment of a counterweight to a rotating shaft in a way that simplifies the attachment.\nCounterweights are often utilized on rotating shafts to balance forces that might be occurring along the shaft as the shaft rotates. As one example, in a scroll compressor, a shaft is driven by an electric motor to drive an orbiting scroll member. The orbiting scroll member is driven to orbit relative to a second scroll member. The counterweight is utilized to balance forces that are created by the orbiting scroll member.\nTypically, counterweights have been attached to the shaft by any one of several attachment methods. As an example, the counterweights have been press-fit or shrunk-fit onto the shafts. Further, staking, swaging or other positive attachment methods have been utilized. These methods have required the counterweight to be formed rather precisely, and further have added the additional costs of the individual attachment methods.\nIt would be desirable to reduce the cost and complexity of attaching the counterweight to the shaft."}
{"text": "1. Field of the Invention\nThe present invention relates to a filter, and more particularly, to a filter including a plurality of LC parallel resonators.\n2. Description of the Related Art\nAs an invention concerning a known filter, a multilayer band pass filter disclosed in Japanese Unexamined Patent Application Publication No. 2011-71921, for example, is known. FIG. 14 is an exploded perspective view of a multilayer band pass filter 500 disclosed in Japanese Unexamined Patent Application Publication No. 2011-71921.\nThe multilayer band pass filter 500 includes dielectric layers 502a through 502g and LC parallel resonators 504 and 516. The dielectric layers 502a through 502g are formed in a rectangular shape and are stacked on each other from top to bottom in this order.\nThe LC parallel resonator 504 includes an inductor electrode 506, via electrodes 508 and 510, a capacitor electrode 512, and a ground electrode 514. The capacitor electrode 512 is disposed on the dielectric layer 502f. The ground electrode 514 is disposed on the dielectric layer 502g. The capacitor electrode 512 and the ground electrode 514 oppose each other with the dielectric layer 502f therebetween so as to form a capacitor.\nThe inductor electrode 506 is a linear conductor disposed on the dielectric layer 502b and extending in the front-and-rear direction. The via electrode 508 passes through the dielectric layers 502b through 502e in a direction in which the dielectric layers 502b through 502e are stacked. The top end of the via electrode 508 is connected to the rear end of the inductor electrode 506. The bottom end of the via electrode 508 is connected to the capacitor electrode 512. The via electrode 510 passes through the dielectric layers 502b through 502f in a direction in which the dielectric layers 502b through 502f are stacked. The top end of the via electrode 510 is connected to the front end of the inductor electrode 506. The bottom end of the via electrode 510 is connected to the ground electrode 514. With this configuration, the inductor electrode 506 and the via electrodes 508 and 510 form an inductor.\nThe LC parallel resonator 516 includes an inductor electrode 518, via electrodes 520 and 522, and capacitor electrodes 514 and 524. The capacitor electrode 524 is disposed on the dielectric layer 502f. The ground electrode 514 and the capacitor electrode 524 oppose each other with the dielectric layer 502f therebetween so as to form a capacitor.\nThe inductor electrode 518 is a linear conductor disposed on the dielectric layer 502b and extending in the front-and-rear direction. The via electrode 520 passes through the dielectric layers 502b through 502e in a direction in which the dielectric layers 502b through 502e are stacked. The top end of the via electrode 520 is connected to the rear end of the inductor electrode 518. The bottom end of the via electrode 520 is connected to the capacitor electrode 524. The via electrode 522 passes through the dielectric layers 502b through 502f in a direction in which the dielectric layers 502b through 502f are stacked. The top end of the via electrode 522 is connected to the front end of the inductor electrode 518. The bottom end of the via electrode 522 is connected to the capacitor electrode 514. With this configuration, the inductor electrode 518 and the via electrodes 520 and 522 form an inductor.\nIn the multilayer band pass filter 500 configured as described above, the two LC parallel resonators 504 and 516 are disposed side by side in the right-and-left direction. With this arrangement, the LC parallel resonators 504 and 516 are electromagnetically coupled to each other so as to form a band pass filter.\nIn the multilayer band pass filter 500 disclosed in Japanese Unexamined Patent Application Publication No. 2011-71921, it is difficult to intensify capacitive coupling between the LC parallel resonators 504 and 516. This will be described more specifically. In the multilayer band pass filter 500, the capacitive coupling between the LC parallel resonators 504 and 516 is adjusted in order to obtain a desired transmission characteristic. If it is desired that the pass bandwidth of the multilayer band pass filter 500 will be increased, intensifying of capacitive coupling between the LC parallel resonators 504 and 516 is effective. For intensifying capacitive coupling between the LC parallel resonators 504 and 516, the distance between the LC parallel resonators 504 and 516 is set to be decreased. Then, the capacitance formed between the via-hole electrodes 508 and 520 is increased, and the capacitance formed between the via-hole electrodes 510 and 522 is increased. As a result, a signal of a lower frequency side than the pass band is more likely to pass between the LC parallel resonators 504 and 516, thereby increasing the pass bandwidth of the multilayer band pass filter 500.\nHowever, if the distance between the LC parallel resonators 504 and 516 is excessively decreased, short-circuiting may occur between the via-hole electrodes 508 and 520 and between the via-hole electrodes 510 and 522. Thus, in the multilayer band pass filter 500, it may be difficult to intensify capacitive coupling between the LC parallel resonators 504 and 516 to satisfy a desired frequency characteristic."}
{"text": "The proliferation of devices that incorporate some form of display, such as televisions, cameras, computer monitors, and portable devices, has become widespread. Despite variations that exist among these devices, they provide a platform to display various kinds of information. While technological advancements towards enhancing the clarity and resolution of display have emerged, and content formats have been standardized, the display quality of various kinds of image data remains problematic."}
{"text": "Biological cells and tissues feature extremely complex three-dimensional (3D) structures spanning a large range of length scale. Living organisms are also very dynamic; molecular processes take place in a timescale covering over 15 orders of magnitudes while the organisms move and metabolize. For example, molecular processes such as protein, DNA, and RNA conformation changes take place in a timescale ranging from 100 fs to 100 s. (For discussion, see, e.g., Whitesides, G. M., Nat. Biotechnol. 21, 1161-1165 (2003); Williams, S. et al., Biochemistry 35, 691-697 (1996); Gilmanshin, R., et al., Proc. Natl. Acad. Sci. U.S.A. 94, 3709-3713 (1997); Callender, R. H., et al., Annual Review of Physical Chemistry 49, 173-202 (1998); Trifonov, A. et al., Journal of Physical Chemistry B 109, 19490-19495 (2005); Cheatham, T. E., Curr. Opin. Struct. Biol. 14, 360-367 (2004); Millar, D. P., Curr. Opin. Struct. Biol. 6, 322-326 (1996); and Brauns, E. B., et al., Physical Review Letters 88 (2002), the disclosures of which are incorporated herein by reference.) The understanding of these processes not only has fundamental biological significance, but also will enable the treatment of a host of human diseases. Essentially, probing these processes requires a “whole-field” technique, which means taking 3D microscopic pictures in a rapid succession with high sensitivity and specificity without scanning, i.e. a four-dimensional (4D) microscope is needed.\nSince holography was invented 60 years ago, an aberration-free, three-dimensional microscope has long been sought. (See, e.g., D. Gabor, “A new microscopic principle”, Nature 161, 777-778 (1948), the disclosure of which is incorporated herein by reference.) In fact, soon after Leith and Upatniek demonstrated the first 3D image by holography, holographic microscopy was proposed. With the recent advent of digital holography, holographic 3D microscopy received renewed and rapidly increasing attention. (E. N. Leith and Upatniek J., J. Opt. Soc. Am. 54, 1295-& (1964); E. N. Leith, et al., J. Opt. Soc. Am. 55, 981-& (1965); U. Schnars and W. Juptner, Appl. Opt. 33, 179-181 (1994); I. Yamaguchi and T. Zhang, Opt. Lett. 22, 1268-1270 (1997); E. Cuche, et al., Appl. Opt. 38, 6994-7001 (1999); F. Dubois, et al., Appl. Opt. 38, 7085-7094 (1999); P. Marquet, et al., Opt. Lett. 30, 468-470 (2005); L. Miccio, et al., Appl. Phys. Lett. 90, 3 (2007); L. Toth and S. A. Collins, Appl. Phys. Lett. 13, 7-& (1968); W. Xu, et al., Opt. Lett. 28, 164-166 (2003); and T. Zhang and I. Yamaguchi, Opt. Lett. 23, 1221-1223 (1998), the disclosure of which are incorporated herein by reference.) Holography is a “whole-field” technique, which means recording all 3D pixels simultaneously in one laser shot without scanning, an extremely valuable asset for biomedical imaging. When combined with modern high-repetition rate pulsed laser and fast imaging devices, a time sequence of consecutive 3D images can be captured, forming a 4D microscope. Indeed, successful 4D holographic imaging has been achieved recently in the context of fluid velocity measurement. (See, e.g., Y. Pu and H. Meng, Appl. Opt. 44, 7697-7708 (2005), the disclosure of which is incorporated herein by reference.)\nDespite the technical advancements, holographic microscopes are not widely deployed in biomedical research because of the lack of specificity. In a microscopic setting with biological samples, holography alone suffers severe background scatterings from irrelevant cell structures. A holographic microscope would capture all scattering entities in the viewing field faithfully, but indiscriminately. On the other hand, the signals of interest (often from small nanostructures like protein molecules) are usually very weak and buried in the strong ambient scatterings from much larger organelles.\nA key to achieve the specificity required for molecular biomedical imaging is to create a contrast between the useful signal and the ambient scatterings. Contrast imaging is routinely achieved through tagging the molecule or nanostructure of interest with fluorescent agents, such as fluorescent dyes, green fluorescent proteins (GFPs), and quantum dots (QDs). By converting the light signal into a different frequency, the unwanted ambient scattering can be easily removed with proper optical filters.\nRecent advances in fluorescence microscopy have profoundly changed how cell and molecular biology is studied in almost every aspect. (See, e.g., Lichtman, J. W. & Conchelo, J. A., Nat. Methods 2, 910-919 (2005); Michalet, X. et al., Annu. Rev. Biophys. Biomolec. Struct. 32, 161-182 (2003); Jares-Erijman, E. A. & Jovin, T. M., Nat. Biotechnol. 21, 1387-1395 (2003); Bastiaens, P. I. H. & Squire, A., Trends Cell Biol. 9, 48-52 (1999); Suhling, K., et al., Photochem. Photobiol. Sci. 4, 13-22 (2005); Chalfie, M., et al., Science 263, 802-805 (1994); Pollok, B. A. & Heim, R., Trends Cell Biol. 9, 57-60 (1999); Jaiswal, J. K., et al., Nat. Methods 1, 73-78 (2004); Michalet, X. et al., Science 307, 538-544 (2005); Zipfel., W. R., et al., Nature Biotechnol. 21, 1368-1376 (2003); Shi, S. H. et al., Science 284, 1811-1816 (1999); Maletic-Savatic, M., et al., Science 283, 1923-1927 (1999); Miyawaki, A., et al., Proc. Natl. Acad. Sci. U.S.A. 96, 2135-2140 (1999); Akerman, M. E., et al., Proc. Natl. Acad. Sci. U.S.A. 99, 12617-12621 (2002); and Alivisatos, P., Nat. Biotechnol. 22, 47-52 (2004), the disclosures of which are incorporated herein by reference.) However, as an incoherent optical process, fluorescence lacks the capability of 3D representation. Should 3D information be required, point-scanning technique has to be performed over the entire volume of interest. The elapsed time for the scanning (i.e. the 3D framing time) is usually in the order of seconds to minutes, which severely constrains the use of this technique. Attempts to extend fluorescence microscopy into three spatial dimensions over time (4D microscopy) are thus incapable of capturing dynamic events. (See, e.g., D. Gerlich and J. Ellenberg, Nat. Cell Biol. 5, S14-S19 (2003), the disclosure of which is incorporated herein by reference.)\nAccordingly, a new holographic imaging methodology and system is needed that allows for dynamic 4D imaging with high contrast and ample spatial and temporal resolution."}
{"text": "Typically, a loudspeaker unit exhibits a low frequency fall-off in its frequency transfer function. In audio systems, it is known to boost low frequencies by having the option of applying a fixed dedicated ‘extra gain’ at low frequencies (‘Loudness’ button). WO9318626A1 describes dynamic adjustment of low frequency gain, where the applied gain is dependent on the content (magnitude) of signal components at low frequencies in the current signal."}
{"text": "Fibreglass as a constructional material enjoys a broad range of characteristics which render it a useful and versatile medium. In particular, fibreglass is readily mouldable into a tough, durable material having highly predictable strength characteristics. Fibreglass has a high strength to weight ratio and good thermal and electrical insulating properties. Fibreglass also has excellent resistance to corrosion and other beneficial properties which render it a particularly attractive material for the construction of ladders and the like.\nHowever, among the restrictions that fibreglass suffers is a limitation applied to its workability by virtue of its lack of ductility. Ductility allows a material to be pressed and pushed into shape. Hence, in the construction of multiple molding constructs of fibreglass, any joints must be either welded, glued or fastened together to effect a reliable joint. Such requirements place a considerable limitation on the type of joining economically available to fibreglass constructs and in particular to small surface area, medium to high stress joints exposed to shear forces and bending about a plurality of axes, for example the rungs of a ladder as they are joined to the stiles. Such joints are difficult to execute in fibreglass as the usual method of pressing ferrules formed in the rungs abutting the stiles cannot be formed on a non-ductile material like fibreglass. The use of glues or welding is limited by the difficulty of injecting sufficient glue into the joint so as to effect maximum bonding of the small surface areas involved."}
{"text": "More than 90% of human deaths resulting from mushroom poisoning are due to peptide toxins found in Amanita species of mushrooms, such as A. phalloides, A. bisporigera, A. ocreata, and A. virosa. Animals, especially dogs, are frequent victims of poisoning by Amanita mushrooms. Recently, two dogs died after eating toxin containing mushrooms in Michigan in the last few months, See Schneider: Mushroom in backyard kills curious puppy, Lansing State Journal, Sep. 30, 2008 (at lansingstatejournal.com.apps/pbcs.dll/article?AID=/20080930/COLUMNISTS09/809300 321.\nHigh concentrations of peptide toxins are found in the above ground mushroom portion (otherwise known as carpophores or fruiting bodies) of the toxin producing Amanita species. These toxins include two major families of compounds called amatoxins (for example, α-amanitin, FIG. 1A) and phallotoxins (for example, phalloidin, phallacidin, FIG. 1B). Both classes of compounds are bicyclic peptides with a Cys-Trp cross-bridge. In general, amatoxins are 8 amino acids in length while phallotoxins are 7 amino acids in length. Although phallotoxins are toxic when injected, phallotoxins do not survive the human intestinal tract and therefore are usually not responsible for deadly mushroom poisonings in humans and animals. On the other hand, amatoxins do survive cooking and remain intact in the intestinal tract where they are absorbed into the body where large doses irreversibly damage the liver. Liver failure due to poisoning by amatoxins can be “cured” only with a liver transplant (Enjalbert et al., (2002) J. Toxicol. Clin. Toxicol. 40:715; herein incorporated by reference).\nThere are an estimated 900-1000 species of Amanita, of which the majority do not produce amatoxins or phallotoxins and some are actually safe for humans to eat (Bas, (1969) Persoonia 5:285; Tulloss et al., (2000) Micologico G. Bresadola, 43:13; Weiβ et al., (1998) Can J. Bot. 76:1170; all of which are herein incorporated by reference). Thus mere ingestion of an Amanita mushroom may not herald the need for the extreme medical treatment necessary to save a patient.\nEven experienced mycologists may not be able to distinguish edible from poisonous mushrooms even with microscopic examination (EMedicine webmd at hypertext transfer protocol site: emedicine.com/ and hypertext transfer protocol site: emedicine.com/emerg/topic874.htm.\nThus physicians and veterinarians need to be able to directly and quickly confirm whether a patient or an animal showing gastrointestinal symptoms of unknown origin, or who has accidentally eaten an unknown mushroom, is in danger of serious illness or death from eating a deadly poisonous mushroom containing amatoxins."}
{"text": "For centuries, silver metal has been known to be an agent capable of killing many different microbial species. It was commonly used to purify drinking solutions or administered to sick individuals before the existence of modern antibiotics. Even after the discovery of penicillin and its descendents, colloidal silver solutions were often used in cases in which troublesome bacteria had become resistant to antibiotics.\nColloidal silver solutions are commercially available today. They are often unstable, however, and have a short shelf life. This is due to the tendency of the silver particles to aggregate and form clusters so large that they are no longer suspended in solution. For this reason, undesirable gelling agents are added to solutions to keep the silver particles suspended by preventing particle aggregation. Another problem of the commercially available solutions is that the majority of the silver content is usually found to be silver ions. This poses a large problem in medical applications where silver ions rapidly combine with ubiquitous chloride to form an insoluble white precipitate.\nNanoparticles have been known to be used as fillers as disclosed in U.S. Pat. No. 6,492,453, as coatings as disclosed in U.S. Ser. No. 2003/0185964 and as foam components as disclosed in U.S. Pat. No. 6,518,324.\nNanoparticle systems are disclosed in U.S. Ser. No. 2002/0150678 as being used in a composition and a method to impart surface modifying benefits to soft and hard surfaces. In particular, this application discloses a soft surface coating for articles such as fabrics and garments.\nInorganic particulates, such as, clays, silicates, and alumina have been widely used in combination with adjunct detergent and laundry compounds to impart some form of antistatic control and/or fabric softening benefit.\nThe present invention relates to composite materials comprising metal loaded onto exfoliated nanoparticles. Such functionalized nanoparticles may be incorporated into soft surface coatings to enhance or modify their bulk physical and performance characteristics. Such soft surface coatings may in turn be used in the preparation of absorbent articles with improved properties. Addition of the coatings to for example the absorbent core of a disposable, absorbent article may help control malodor formation and increase absorbency.\nIn one embodiment, the metal is silver and the nanoparticle comprises a nanoclay. Silver ion is reduced to its neutral metal state (Ag0) and loaded onto the nanoclay. Silver-coated nanoclays in particular have excellent antimicrobial properties, and represent a less expensive alternative to the use of colloidal silver solutions. Such nanoparticles made according to the invention are stable and use less silver metal to generate the same surface area as solid silver particles, making them more cost efficient."}
{"text": "This invention relates generally to an improved speaker apparatus and in particular to a speaker apparatus having a main speaker and at least one tweeter speaker in overlying relation thereto.\nIn the past it has been known to assemble a main speaker and a tweeter speaker in a combined assembly by the use of sheet metal parts secured by rivets, nuts and bolts and washers. A problem associated with the sheet metal configuration has been blockage of the main speaker. Furthermore, the fastening parts have been known to vibrate and loosen and thereby cause a disturbance which detracts from the quality of the speaker assembly."}
{"text": "This invention is an improvement of a mechanical locking device specified in U.S. patent application Ser. No. 08/988,827, which is an improvement over U.S. patent application Ser. No. 08/782,611, then allowed Apr. 16, 1989.\nThe allowed Patent Application '611 provides a gearshift lever locking device which comprises basically a relatively bulky main body housing most of the mechanical parts including sleeve, locking mechanism, hinge and others in a single unit. The bulky size of this main body requires a very large clearance in the gearshift lever which hindrance the applications of this anti-theft device in automobiles of relatively small gearshift lever clearances. The improvement made by Patent Application '827 moves the locking mechanism to the latching member and compresses the main body into a relatively small and simple sleeve and hinge structure. This modified anti-theft device can then be fitted into gearshift levers of very small clearance, resolving the first and foremost problem of adaptation and application of this type of anti-theft devices.\nDespite the above improvement, however, the problems of adaptation are not limited to the clearance factor alone. Some of the others are: size of the lever, tilt angle between the lever and the slot, and length and width of such slot. That is: there are different sizes such as the diameters and shapes for the gearshift or control levers; the tilt angle between the lever and the slot may be different from 90 degrees; there may be different widths and lengths for the slot. A universally adaptable device should have mechanisms adjustable to all these factors. The present invention is developed to meet with this requirement.\nVery briefly, I have found that the adaptation of an interlocking sleeve set with U-shaped channel of various diameters will take care of the size factor of the control lever. By means of a yoke structure slidably received in a latching member and moving the locking mechanism including the lock tumbler away from the base member, the tilt angle and slot length factors may be covered. Lastly, the introduction of a jacketed sheath structure will resolve the variations in the width of the slot. In this way, a new and versatile anti-theft device in the kit form is therefore resulted. I have also found that such a novel set of locking device may not only be applicable to gearshift lever of an automatic transmission for an automobile but also may be easily adaptable to gearshift lever of a conventional or standard shift automobile or even control levers of many machinery's operating in a slot structure. I therefore name this device as a control lever locking kit."}
{"text": "Heretofore various forms of locking type bolt and nut constructions have been designed. Various forms of these constructions are disclosed in U.S. Pat. Nos. 113,557, 3,176,746, 3,339,389, 3,477,337, and 3,517,717. However, for various reasons these previously known locking bolt and nut threaded connections have proven inefficient, at least under some working conditions."}
{"text": "1. Field of the Invention\nThe invention relates to a vehicular drive control apparatus and a method for controlling a vehicular drive apparatus. More particularly, the invention relates to a vehicular drive control apparatus which is provided with a first transmitting portion and a second transmitting portion, and reduces shift shock when a shift is performed in the first transmitting portion at the same time that a shift is performed in the second transmitting portion. The invention also relates to a method for controlling a vehicular drive apparatus provided with a first transmitting portion and a second transmitting portion, which reduces shift shock when a shift is performed in the first transmitting portion at the same time that a shift is performed in the second transmitting portion.\n2. Description of the Related Art\nJapanese Patent Application Publication No. 2005-337491 (JP-A-2005-337491) describes a vehicular drive apparatus in which an internal combustion engine, a first transmitting portion, a second transmitting portion, and driving wheels are connected in series. According to the clutch and brake application chart shown in FIG. 2 of JP-A-2005-337491, in a control apparatus of that vehicular drive apparatus there is no shift pattern in which the first transmitting portion and the second transmitting portion are shifted simultaneously in a normal shift operation. For example, in FIG. 2 of JP-A-2005-337491, during a shift from 4th gear to 5th gear, a shift is performed in the first transmitting portion but not in the second transmitting portion at that time. However, although not common knowledge, a shift pattern in which the first transmitting portion and the second transmitting portion are shifted simultaneously is conceivable depending on the structure of the planetary gear sets of the first and second transmitting portions.\nTypically, when a shift is performed in an automatic transmission, shift shock occurs at the end of the shift, i.e., at the end of an inertia phase when the input rotation speed of the automatic transmission changes. Therefore, when a control apparatus of a vehicular drive apparatus has a shift pattern in which the first transmitting portion and the second transmitting portion are shifted at the same time, and the first and second transmitting portions are shifted at the same time and those shifts end at the same time, it is likely that an occupant will feel significant shift shock such that comfort is lost. However, the control apparatus of the vehicular drive apparatus described in JP-A-2005-337491 does not have a shift pattern in which the first transmitting portion and the second transmitting portion are shifted at the same time. Therefore there is no mention in JP-A-2005-337491 of control to reduce the potential discomfort due to significant shift shock being felt by an occupant as a result of shifts in the first and second transmitting portions ending at the same time."}
{"text": "The present invention relates to a control device for wearable electronics devices and systems, that is to say, to an electronic device configured such as to be incorporates into conventional clothing, and designed so as to be comfortable for the user to wear. This comfort may arise through the avoidance of flat, rigid surfaces, but preferably comes from the use, so far as possible, of flexible parts conformable to the human body.\nExamples of wearable electronics are given in the commonly-assigned UK patent application number 9927842.6 filed 26th Nov. 1999 and entitled xe2x80x9cImproved Fabric Antennaxe2x80x9d, and U.S. Pat. Nos. 5,798,907 and 5,912,653.\nFor such electronic devices, control by conventional rotary knobs, slider potentiometers, and rocker switches is often inconvenient, and especially difficult when the wearable electronics device is intended for use by those having limited manual dexterity, whether through age or infinity, or through the device being incorporated into outdoor winter garments, when the wearer is likely to be wearing gloves or mittens.\nA particularly useful feature in wearable electronics is the ability to adjust performance in response to detected movements of the user. To this end, it is known to provide a fabric which has an electrical resistance which varies as a fabric is deformed: U.S. Pat. No. 4,715,235 (assigned to Asahi Kasei Kogyo Kabushiki Kaisha) discloses a deformation sensitive electroconductive knitted or woven fabric. In the described fabric construction, the electroconductivity is changed when stretch or compression is applied to the fabric, and the provision of electrodes at two or more positions on the fabric allows the degree of deformation to be sensed electrically. The disclosure of U.S. Pat. No. 4,715,235 relates mainly to the construction of the fabric; applications mentioned include a device for detecting stretch or flex in a joint of a human body or human respiration; as a finger switch mounted to lie along a finger portion of a glove and operable by flexing the finger; as a touch sensor on a window frame or hand rail to detect criminal entry; and as an industrial sensor to control a robot.\nIt is an object of the present invention to provide a control device for wearable electronics which is simple and intuitive in use, and which can be used by an impaired or gloved hand.\nIn accordance with the present invention there is provided a manual control for an electronic device which has control means to vary a variable of the device having at least one cord electrically connected to the control means, the cord having an electrical property which varies in accordance with the tension applied to the cord.\nAlso according to the invention, there is provided a garment for use with an electronic device which has a control means to vary a variable of the device, the garment comprising means to support the electronic device, at least one cord having a first end connected electrically to the control means, the cord having an electrical property which varies in accordance with the tension applied to a second end of the cord, and the control means being arranged to sense the electrical property of the cord and to vary the variable of the electrical device accordingly.\nFurther features and advantages of the present invention are defined in the attached claims or will become apparent from reading of the following description of embodiments of the present invention."}
{"text": "This invention relates to a method for protecting plants from attack for extended time periods by soilborne plant disease organisms which comprises contacting seeds or plant parts with a systemic plant protectant corresponding to the formula ##STR2## wherein X represents chloro, fluoro or bromo; Z represents trichloromethyl, dichloromethyl or dichlorofluoromethyl and R represents lower alkyl of 1 to 4 carbon atoms or phenyl.\nIn U.S. Pat. No. 3,244,722, issued Apr. 5, 1966, there are described and claimed, among other related compounds, those corresponding to the formula: ##STR3## wherein R is alkyl of 1 to 18 carbon atoms.\nExemplary compounds listed in this patent include 2-chloro-4-methoxy-6-(trichloromethyl)pyridine, 2-chloro-6-methoxy-4-(trichloromethyl)pyridine, 5-chloro-2-methoxy-4-(trichloromethyl)pyridine and 3-chloro-2-methoxy-4-(trichloromethyl)pyridine. As reported in this patent, various compounds disclosed therein are useful as herbicides; various other compounds are useful in the control of pest fish and aquatic insects; and other compounds are taught to be useful as insecticides and anthelmintic agents for warm-blooded animals.\nIt has now been found that within the broad genus of pyridine compounds taught in U.S. Pat. No. 3,244,722, the specific compounds having a chlorine, fluorine or bromine atom in the 6-ring position, a dichloromethyl, trichloromethyl or dichlorofluoromethyl group in the Y-ring position and an alkoxy or phenoxy group in the 2-ring position surprisingly and unexpectedly exhibit superior and unique properties in being able to systemically protect plants for extended periods of time from attack by soil-borne plant diseases. This kind of activity is not exhibited by isomers or analogs of compounds employed in the herein disclosed method. The outstanding activity of these specific compounds is particularly unexpected since these compounds are generally not effective at the same application rates against the same organisms outside of the plant itself, i.e., when the compounds are used as a pre-plant soil sterilant or fumigant.\nIt is an object of the present invention to provide an improved method for systemically protecting plants from attack, for extended periods of time, by soil-borne plant disease organisms.\nIt is a further object of the present invention to provide a method whereby long-lived activity in plants against soil-borne plant disease organisms is achieved, with toxicants which are generally not effective in controlling these same organisms outside of the plants.\nOne of the advantages of the present method is that by the mode of action of the active toxicant, plant diseases can be eliminated from infected plants and non-infected plants can be protected from attack.\nThe present method also offers a practical advantage in that there is no need to employ the additional time and labor required by conventional preplant sterilization with soil fumigants.\nA further practical advantage of the present method is that the active compounds are used in amounts which are the equivalent of ounces to several pounds of the active ingredient on a per acre basis as against the conventional soil fertilization practices which require many pounds and hundreds of pounds of active material per acre."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit for discriminating trimming failure.\n2. Description of the Related Art\nAn example of conventional circuit having trimming structure is shown in FIG. 8. In FIG. 8, the resistance between a terminal 811 and a negative power supply VSS corresponds to a resistance 841 by cutting a fuse 821. Further, if the fuse 821 is not cut, the resistance between the terminal 811 and the negative power supply VSS corresponds to a combined resistance of the resistance 841 and a resistance 842. A voltage of the terminal 811 of the circuit in FIG. 8 is controlled by the resistance ratio to a resistance 843.\nIn the conventional trimming structure, a non-conforming article is generally produced if trimming is not performed correctly. However, there is a possibility that a conforming article is nonetheless provided without performing trimming depending on a range of an efficiency voltage of a test. In this case, even if a laser is radiated to a position away from the fuse location so that trimming fails to be performed because of positional deviation of trimming alignment, or the like, a conforming article may be provided at the time of the test. In addition, there is a fear that the radiation of a laser to the wrong position affects long-term reliability of the device, and thus, the non-conforming article has to be correctly discriminated.\nIn order to solve the above-described problems, an object of the present invention is to provide a semiconductor integrated circuit having a structure in which a dummy fuse is used to discriminate failure in accordance with an output thereof in the case of a trimming error. The semiconductor integrated circuit as structured above has a strong point in that success or failure of trimming can be discriminated in accordance with whether the output is fixed to VSS or VDD or not. Further, such a structure is adopted in which, if trimming is not performed, short circuit is established between positive and negative power supply voltages, whereby the success or failure of trimming can be discriminated by examining a consumption current."}
{"text": "1. Field of the Invention\nThis invention relates to a method and apparatus for summing the products of samples of analog signals, to discrete analog signal processing (DASP) systems and, more particularly, to systems employing CCD devices configured for implementing a real-time, fully-analog correlator or convolver.\n2. State of the Prior Art\nMatrix multiplication of analog signals and, in particular, analog signal correlation, for example, has been performed by such techniques as converting the analog signals to digital signals and performing digital multiplication or correlation. However, such A/D conversion techniques require large amounts of hardware which uses relatively large amounts of space and decreases system reliability.\nIt is desirable, therefore, to devise a method and apparatus for analog signal matrix multiplication or correlation which does not include A/D conversion.\nAn application, concurrently filed herewith, and entitled \"CMOS Analog Multiplier for CCD Signal Processing\" of Lampe, et al., is of interest for its disclosure of a balanced differential conductance-type MOS transistor multiplier using complementary MOS transistors.\nAlso of interest is U.S. Pat. No. 4,035,329 by Lampe, et al., assigned to the common assignee which discloses an extended correlated doubling sampling system which corrects for errors in components of the system.\nApplication Ser. No. 625,696, filed Oct. 29, 1975, entitled \"A Dual-CCD Real Time, Fully-Analog Correlator\", by Lampe, et al., is of interest for its disclosure of a general-purpose fully analog correlator including two CCD channels for processing samples of analog signals and MOSFET multipliers for providing products of the signal samples."}
{"text": "1. Field of the Invention\nThe present invention relates to keyboard operated musical instruments. More particularly, the present invention relates to interfacing an acoustic or electrically amplified piano to a digital computer, storage device, or music synthesizer.\n2. Description of the Prior Art\nSince before the turn of the century, efforts have been made to capture information from acoustic keyboard instruments as the instruments are played. Early devices of this type were predominantly player pianos wherein a performance could be captured in the form of holes on paper that controlled a pneumatic piano control system. The point of such devices was to make one recording on a paper roll of the performance. Thereafter, the performance could be replayed as desired on a player piano by use of a copy of the original paper roll.\nOne more recent approach to capturing keyboard information during a performance employs a switch matrix including a single matrix row coupled to a bus extending underneath a piano keyboard and having a number of matrix columns corresponding to the number of keys on the instrument. The switch matrix includes a contact on each key operable to contact the matrix bus when the key is operated and thus indicate a switch closure. Such device is shown in U.S. Pat. No. 3,604,299, issued to Englund, Sept. 14, 1971.\nThe installation of such device is time consuming and, therefore, expensive. Once installed, it is likely that the action of the piano may be affected. Additionally, the open contact and bus system is subject to performance degradation and lack of reliability as the result of oxidation and contact wear. A similar arrangement is provided in U.S. Pat. No. 4,023,456, issued to Groeschel on May 17, 1977.\nU.S. Pat. No. 4,104,949, issued to Clark on Aug. 8, 1978, provides a series of keys located on the piano above the keyboard and including sensors responsive to keyboard operation. The location of such device above the keyboard eliminates some of the playing area on the keyboard available to the musician and, thus, interferes with normal keyboard operation.\nThe most recent approach to capturing keyboard information in a mechanical or electrical manner is shown in U.S. Pat. No. 4,307,648, issued to Stahnke on Dec. 29, 1971 and in U.S. Pat. No. 4,351,221, issued to Starnes et al. on Sept. 28, 1982. Both such devices use an optical switch including a light source and photo cell. Key operation results in movement of a key or hammer within the piano mechanism that produces a corresponding movement of a flag attached thereto. As a result, a beam of light between the light source and photo cell is broken by interposition of the flag, and information is produced corresponding to key operation.\nSuch devices are installed to the piano action during a time consuming and, therefore, expensive installation procedure. Because there is continual movement of the flag, there is a likelihood that the flag will become dislodged or slip out of alignment and not operate the optical switch properly. Thus, such systems generally require considerable maintenance. Furthermore, the addition of a flag to the piano action might, in some circumstances, affect the action of the piano and, thus, alter or degrade the feel and playing quality of the piano.\nAn additional problem in prior art keyboard interfaces is that of cabling. A standard piano keyboard has eighty-eight keys. Known keyboard interfaces, therefore, require a cable assembly including eighty-nine wires (eighty-eight key wires and one return or ground wire). Although there have been some attempts to reduce this number of wires through various decoding and latching schemes, generally, significant problems are introduced by the prior art devices when routing cables to transfer electronic signals in response to keyboard operation.\nElectronic musical instruments include the well-known electronic music synthesizer, a device that is capable of imitating many known electronic and acoustic instruments, as well as producing its own unique sounds. Electronic music synthesizers are operated by keyboards, or slide controllers or both. One disadvantage of an electronic music synthesizer is that the keyboard generally provided is of a size (that is, number of keys), and has an action or feel unlike that of acoustic pianos. Many musicians have studied on acoustic pianos. A musician who has become accustomed to the feel of a particular keyboard tends to prefer that particular keyboard over other keyboards having a different feel.\nAnother disadvantage of synthesizers is that they introduce an additional keyboard into a playing configuration. For example, during a performance, a musician who plays both an acoustic keyboard instrument and an electronic music synthesizer has to physically shift between the two. Such additional keyboards take up space in the club or theater where the performance is to be held and provide extra equipment to be transported from location to location in accordance with the musician's itinerary."}
{"text": "1. Field of the Invention\nThis invention relates to a driving method and device for generating activating signals that serve to activate scan lines of a display panel.\n2. Description of the Related Art\nLiquid crystal displays (LCD) are well known in the art and have largely replaced cathode-ray tubes (CRT) in view of their thin profile and lightweight characteristics. The LCD may be of the twisted nematic (TN) type, the super-twisted nematic (STN) type, or the thin film transistor (TFT) type. The TFT type LCD is relatively expensive to manufacture when compared to the TN and STN type LCD. Therefore, the TN and STN type LCD are often used for mobile phones and digital cameras. FIG. 1 illustrates a conventional STN type LCD that includes a display panel 92 and a driving device 91. The display panel 92 has a plurality of scan lines (Com1 to ComN). The driving device 91 is connected to the scan lines (Cora1 to ComN) of the display panel 92, and is 2 operable so as to generate activating signals that serve to activate the scan lines (Com1 to ComN) of the display panel 92 in a progressive scanning manner.\nThe aforementioned conventional STN type LCD is disadvantageous in that, with further reference to FIG. 2, the driving device 91 thereof generates the activating signals such that a time point at which a pulse duration of a preceding one of the activating signals, e.g., the activating signal (Si), in a consecutive pair of the activating signals ends and a time point at which a pulse duration of a succeeding one of the activating signals, i.e., the activating signal (S2), in the consecutive pair starts occur simultaneously, e.g., at a time point (t32). However, due to non-ideal effects, e.g., uneven loading effects, of the scan lines (Com1 to ComN), the activating signals (S1, S2) are attenuated, and as a consequence, the time point (t33) at which the pulse duration of the activating signal (S1′) ends may occur after the time point (t32) at which the pulse duration of the activating signal (S2′) starts. That is, the activating signal (S2′)begins to activate one of the scan lines (Com1 to ComN), e. g. , the scan line (Com2), in an adjacent pair of the scan lines (Com1 to ComN) before the activating signal (S1′) has completely deactivated the other of the scan lines (Com1 to ComN), e.g., the scan line (Com1), in the adjacent pair. This easily arises in crosstalk between he scan lines (Com1, Com2) and results in a “shadow” appearing on the display panel 92 of the conventional STN type LCD.\nFurthermore, since the scan lines (Com1 to ComN) have different lengths and therefore different non-ideal effects, the scan lines (Com1 to ComN) attenuate the activating signals at different degrees. In particular, the scan line (Com1), which has the shortest length, attenuates the activating signal (Si) the least, whereas the scan line (ComN), which has the longest length, attenuates the activating signal (Sn) the most.\nTo solve the above-mentioned problems, it has been proposed to gradually increase the widths of the scan lines (Com1 to ComN) from the shortest scan line (Cora1) to the longest scan line (ComN). This approach, however, results in a larger area, e.g., of a circuit board (not shown), occupied by the scan lines (Com1 to ComN). Such problem becomes even more serious for display panels having a large number of scan lines. In addition, since the degree at which the width of the scan line (Corn] to ComN) is increased is determined based on a particular type of the display panel 92, once the scan lines (Com1 to ComN) are formed, e.g., on the circuit board, the widths thereof are not capable of being altered to suit a different type of display panel."}
{"text": "The disclosures of all publications, patents, patent application publications and books referred to in this application are hereby incorporated by reference in their entirety into the subject application to more fully describe the art to which the subject invention pertains.\nThe development of safe and effective therapies for treating acute and chronic wounds is an issue currently of great interest to clinical scientists and industry, alike. Wound healing is an intricate, multi-stage process that relies heavily on the delivery of new cells to the wound zone. Two key elements of the wound healing response are fibroplasia and epithelialization when fibroblasts and epithelial cells, respectively, enter the wound to form a protective barrier from the external environment. This is stimulated by cell proliferation and migration from the wound edge. The identification of agents that increase the rate at which cells invade and close a wound would represent a major advance in wound healing therapeutics. Ideally, this would be a topically applied agent that stimulates the proliferation and migration of fibroblasts and wound edge epithelial cells.\nThe present invention addresses this need and identifies a novel target in promoting wound healing and provides therapies and assays based thereon."}
{"text": "In an image forming apparatus using electrophotographic methods, toner may be supplied to an electrostatic latent image formed on an image receptor so as to form a visual toner image. The visual toner image may be transferred onto a sheet, and the image may be fused onto the sheet using a fusing device that uses heat or pressure. The fusing device may use a fusing nip, which may be formed by engaging a heating roller or a heating belt with a pressuring roller. A sheet may be passed through the fusing nip so that toner may be melted and fused onto the sheet. Alternatively, the fusing device may use light from a light source, which is directly emitted to a sheet so that toner may be melted and fused onto the sheet.\nThe fusing performance of the toner may be closely related to the heat absorbance of the toner. For example, when the toner is fused using a light source as a heater, the absorbance of the toner with respect to optical energy affects the fusing performance of the toner. That is, as the optical absorbance of the toner increases, the warm-up time may be reduced, and power consumption of the image forming apparatus may also be reduced.\nIn order to increase the optical absorbance of the toner, a light absorber may be added to the toner. However, a general light absorber absorbs some of the light having wavelengths in the visible ray range, thereby adversely affecting the reproducibility of color of the toner. Thus, there remains a need in the art for new and improved electrographic toners having improved optical properties."}
{"text": "This invention relates to the field of optics and particularly to materials for use as a nonlinear element in nonlinear optical devices.\nRecent advances in nonlinear optics (NLO) have created a new frontier for applied optics. This technology requires nonlinear optical materials, i.e., materials which alter the frequency of laser light, materials which have an index of refraction that varies with light intensity or with applied electrical field, or (in the case of photorefractive nonlinear optical materials) materials which have a local index of refraction that is changed by the spatial variation of light intensity.\nDevices which have been conceived utilizing NLO materials include parametric amplifiers, oscillators, second harmonic generators, and modulators. Such devices provide, inter alia, second harmonic generation, phase-conjugate navigation, and laser beam combining which employs optical phase conjugation, as well as laser beam spatial (and spectral) mode cleanup, wavelength agile rejection filtering, laser radar, image correlation and enhancement, communications, optical data storage, and optical computing. In particular, there are readily identifiable systems applications which need lower-cost, higher-response, higher-average-power materials for optical parametric amplifier operation and second harmonic generation throughout the blue/near-ultraviolet spectral region. Bierlein, et al., U.S. Pat. No. 3,949,323, for example, depicts, in FIG. 1, the use of a nonlinear optical material to generate second harmonic radiation and, in FIG. 2, the use of such a material to modulate the intensity of a beam of polarized light. Many of these concepts rely on real-time holographic effects which can be induced in nonlinear materials. Although several of these innovations are at preliminary stages of design, such concepts reveal the power of nonlinear optics in terms of performance breakthroughs in electro-optic device technology that a decade ago would have been thought impossible.\nThe manufacturability of nonlinear optical devices, however, is limited by the performance of currently available NLO materials, which are primarily inorganic oxide crystals such as potassium niobate, potassium dihydrogen phosphate, lithium niobate, strontium barium niobate, and barium titanate. These inorganic oxide crystals exhibit shortcomings as nonlinear materials due to difficulty of synthesis, lack of optical quality, and slow electro-optic response times.\nAlthough there are many organic materials (for example, urea and 2-methyl-4-nitroaniline) which exhibit very high second harmonic generation responses in powder form, there have been no practical device-quality organic NLO crystals available in the art. Typical organic crystals lack the mechanical robustness, dimensional and thermal properties, and optical quality required for fabricating optical devices.\nA category of \"semi-organic materials\", or metal-organic coordination complexes, has been reported to show promise in the development of nonlinear optical materials. See, for example, Wenbo et al., A new organometallic nonlinear optical material-triallylthiourea mercury bromide (ATMB) crystal: growth and structure, Journal of Crystal Growth, Volume 133, Page 71 (1993) and Dong, et al., A New Aromatic Organometallic Nonlinear Optical Crystal: [Bis-4-nitropyridine-N-oxide Cadmium Chloride], Materials Research Bulletin, Volume 29, Page 73 (1994). This materials classification encompasses a very large number of semi-organic ionic crystals in which relatively large, polarizable organic molecules are incorporated into a host inorganic ionic lattice.\nAdvanced nonlinear optical (NLO) applications such as high speed optical phase conjugation, parametric amplification, and laser hardening require NLO crystals having high speed and high efficiency. Devices for these applications require new and improved NLO crystals in order to achieve their potential."}
{"text": "1. Field of the Invention\nThe present invention relates to a focus controller and an optical apparatus using the same.\n2. Description of the Related Art\nConventionally, Japanese Patent Laid-Open Nos. 05-150154 and 05-064056 disclose mountain climbing driving as an autofocus (“AF”) for an optical apparatus, such as a video camera or a TV camera. The mountain climbing driving is a method that detects a definition of a photography screen from a high frequency component in a video signal obtained from an image pickup lens, and controls a position of a focus lens in an optical-axis direction so that the definition can become maximum. The conventional position control finely reciprocates (wobbles) the focus lens in the optical-axis direction at a constant period through an open loop control using a stepping motor.\nHowever, the stepping motor loses synchronization between a control pulse and a motor's rotation in the high-speed rotation or with an excessively large load, (stepping out). In order to solve this problem, Japanese Patent Laid-Open No. 10-150798 proposes use of a blushless DC motor instead of the stepping motor. Since the blushless DC motor detects a rotor's position by using a sensor, it has an advantage in maintaining synchronization between the instruction signal and the rotor's position and in preventing stepping out.\nNevertheless, the open loop control using the blushless DC motor cannot control a lens's position, and has a problem in that the wobbling amplitude becomes inaccurate. A closed loop is needed for accurate positioning control with the blushless DC motor, but thereby causes a problem in which a circuit becomes complex."}
{"text": "Recently, multifunction devices produced by adding a number of functions to a digital copier serving as a base have been put to practical use. The multifunction devices have capabilities to conduct facsimile communications using scanner and printer functions inherent in the digital copier, convert code data sent from a computer into bitmap data and print out the resulting bitmap data on the printer, transmit images read by the scanner to other devices on a network, and print images delivered via the network on the printer.\nAlso, control panels of the multifunction devices have become larger in size and higher in resolution, making it possible not only to configure settings for the scanner and printer functions, but also to operate remote apparatus on a network and monitor equipment condition using their operation screens or display or print web content using a built-in browser. Incidentally, patent documents related to the present invention include Japanese Patent Laid-Open No. 11-134125.\nHowever, with built-in web-browser functions of a multifunction device such as those described above or a personal computer, results of web site searches by means of a search engine may vary with the database contents of the search engine even if the same keywords are used.\nOn the WWW, there are always web sites which appear newly, web sites which disappear, and web sites which change their contents. Data in databases of search engines are generally collected by crawling around the WWW automatically. The data are constituted of snapshots taken during the automatic crawling and the database contents can often vary depending on time. Thus, when a user performs a search using search criteria on a search engine, web sites which match the search criteria can vary with the time of the search.\nThus, by analyzing variations in search results with time, it is possible to tell sites which have appeared newly. However, search results are generally listed in an order irrespective of the time of appearance of web sites and if a user tries to analyze variations in search results (difference information) manually, checkups involve a great burden. Also, the user must check URLs one after another visually by comparing them between search results produced at two different times. Although recent multifunction devices and other information processing apparatus are equipped with relatively high-performance screens, they are insufficient in terms of screen size, viewability, or operability, increasing the burden on the user with increases in the volume of search results."}
{"text": "Heretofore, rosins and rosin esters being the derivatives thereof have been used variously as tackifiers for hot-melt adhesives and pressure-sensitive adhesives, modifiers for rubbers and various plastics, base materials for chewing gum, and raw materials such as resins for traffic paints, sizing agents for paper making, emulsifiers for synthetic rubbers, resins for ink and resins for coating, and the like. However, appearance of the above-mentioned rosin compounds was colored in yellow or yellowish brown, and the compounds were not satisfactory from the viewpoints of odor, thermal stability and weatherability (hereinafter referred to thermal stability and weatherability as stability), and the like.\nTherefore, in order to dissolve the above-mentioned problems in the rosin compounds, a disproportionated rosin having an improved thermal stability which is prepared by disproportionating an unpurified rosin, a hydrogenated rosin having a somewhat improved stability which is prepared by hydrogenating an unpurified rosin, or their ester compounds which are prepared by esterifying the above-mentioned disproportionated rosin or hydrogenated rosin are commercially available. However, each and every one of these rosin compounds does not have satisfactory color tone, stability, and the like.\nJapanese Examined Patent Publication No. 33771/1970 and Japanese Examined Patent Publication No. 20599/1974 disclose a method for disproportionating rosins or rosin compounds with a specific organic sulfur compound. However, the disproportionated rosin esters or the disproportionated rosins which are obtained according to these methods are not satisfactory from the viewpoints of color tone, stability and odor.\nJapanese Unexamined Patent Publication No. 9605/1980 discloses a process for preparing rosin esters having a higher stability, comprising purifying a disproportionated rosin to remove high molecular weight materials, unsaponifiable materials and the like which are contained in the raw rosin and successively esterifying the purified disproportionated rosin with an alcohol. The rosin ester obtained according to the process is relatively superior in light color and stability to the known rosin esters. However, the obtained rosin ester is colored in the above-mentioned esterifying step and does not have a satisfactory heat coloration resistance. Accordingly, the obtained rosin ester should be improved upon the properties mentioned above.\nAlso, Japanese Unexamined Patent Publication No. 186783/1988 discloses a process for preparing a colorless rosin ester having a high stability. Japanese Unexamined Patent Publication No. 85265/1989 discloses a process for preparing a colorless rosin having a high stability. However, a hydrogenating step is indispensable to these processes. In the processes, an advanced high-pressure-proof apparatus is necessarily used, or a large quantity of hydrogen is consumed. The cost of preparation is remarkably increased thereby. Therefore, these processes are disadvantageous.\nFurther, Japanese Unexamined Patent Publication No. 230072/1984 discloses a process for preparing a light-colored rosin ester having a high stability, comprising esterifying a purified rosin by distillation with an alcohol in the presence of a specific organic sulfur compound having abilities of both disproportionating and light-coloring. However, the obtained rosin ester does not have satisfactory color tone and stability and there is a problem that the obtained rosin ester has strong sulfuric odor which is generated from the organic sulfur compound when the rosin ester is heated.\nAs mentioned above, any conventional rosin esters, rosins, disproportionated rosins and hydrogenated rosins have not yet been satisfactory in odor, color tone and stability concurrently. In addition, the above-mentioned rosin compounds could not have competed at all in the above-mentioned properties with the hydrogenated petroleum resins, which are employed in the same uses as the rosin compounds. Therefore, it has been desired to find a rosin ester and a rosin which are colorless and have a high stability.\nIn consideration of the above-mentioned circumstances, the present inventors have studied to develop novel processes for preparing a rosin ester and a rosin of which various properties such as color tone, odor and stability have been further improved, in comparison with those of conventional rosin esters and rosins, which enable to cheaply provide such a rosin compound."}
{"text": "Field of the Invention\nThe present invention relates to an induction assembly and system for a supercharged internal combustion engine. More particularly, the present invention relates to an assembly and a system having a monolithic unitary cast housing securing a super charger assembly and three integrated intercooler assemblies allowing for both a variable rate convective heat transfer via air cooling of an exterior of the monolithic cast housing and an air cooling of pressurized discharge air discharged from the super charger assembly for combustion.\nDescription of the Related Art\nIndustrial applications of induction assemblies for supercharged internal combustion engines have included a plurality of complex air inlet runners and cylinder head attachments.\nIn one related matter in U.S. Pat. No. 6,029,637 (Prior), the entire contents of which is incorporated herein by reference, an induction assembly is provided with extended induction housing intakes related to a super charger arrangement. Here, the entrance of the intake runners is long extended providing great inconvenience in access, enhanced costs, and greater difficulty in disassembly and maintenance. Additionally as a detriment, the assembly cannot be provided in a single compact monolithic manner due to the long extended intakes and the requirement for differently shaped cylinder heads and accesses geometries.\nAlso in a similar detrimental arrangement is U.S. Pat. No. 7,426,921 (Billings, et al.), the entire contents of which is incorporated herein by reference, wherein a super charger arrangement is provided with a separate-piece-kit type of two-part air inlet casings that are bolted to a side of a rotor casing in combination with a top rotor casing cover member. The arrangement provided is substantially complex and is also weakened by positioning the separate two-part air inlet casings with bending moments being secured to the rotor casing only through a separate cover member. Several detriments to air flow, rigidity, and sound attenuation. Additionally, thermal release from pressurized air is greatly hampered limiting power gain substantively.\nAccordingly, there is a need for an improved induction assembly and system for a supercharged internal combustion with enhanced process efficiencies and thermal release."}
{"text": "The present invention relates to a method of and apparatus for protecting a rotating machine such as a centrifuge including a rotating body or rotor rotating at high speed.\nHigh-speed rotating machines such as centrifuges tend to produce a self-excited vibration when an excess clearance exists between rotating parts. Furthermore, a rotor of such machines is apt to be in an imbalanced state when a sample is inaccurately set by a user. Under such condition, the rotor of the machine oscillates at a large vibration amplitude and bearings for supporting the rotor are subjected to increased radial loads exerted thereon. To prevent damage or breakage of the bearings, it has been customary to detect the imbalance of the rotor, in terms of the amplitude of vibration of the rotor or a drive shaft connected to the rotor, so that rotation of the rotor is stopped when the detected amplitude of vibration exceeds a predetermined allowable value A.sub.0 as at a point e shown in FIG. 3 of the accompanying drawings. The conventional practice is not satisfactory if the allowable value is set to a relatively low level because the detection of the amplitude of vibration of the rotor is greatly affected by external disturbances. This means that the detected vibration amplitude readily exceeds the low allowable value under the influence of the external disturbance, which causes frequent stopping of the rotor. On the other hand, an allowable value A.sub.0 set at a high level permits the rotor to continue rotating while oscillating at an amplitude of vibration f which is slightly smaller than the allowable range A.sub.0. Under such mode of operation, the bearing life is substantially reduced."}
{"text": "The present invention relates to electroless nickel plating, more particularly to the electroless nickel plating of surfaces such as copper or fused tungsten, and still more particularly to the activation of, e.g., copper or fused tungsten surfaces for receipt of electroless nickel.\nIn many applications of the fabrication of printed circuits having conductive copper circuity paths and areas, as well as in applications of fused tungsten circuitry on ceramic substrates, electroless nickel is desired to be plated over the copper or over the fused tungsten. In printed circuits, for example, electroless nickel plating over copper is often used to form a diffusion barrier for a subsequently-applied gold plating, while in tungsten circuitry electroless nickel plating is used, e.g., to provide more readily solderable surfaces.\nA wide variety of electroless nickel plating baths are available, comprised generally of aqueous solutions containing a source of nickel ions, a reducing agent for the nickel and a complexing agent, capable of operating over predetermined ranges of pH, with the most common such baths being based upon hypophosphite reducing agents. In order to electrolessly plate nickel from such baths onto copper surfaces (e.g., in printed circuit boards), it is necessary to first activate the copper surfaces. The most common means for effecting this activation is by contact of the copper surfaces with species catalytic to deposition, typically colloidal sols or solutions of palladium-tin (see, e.g., the one-step catalyst compositions of U.S. Pat. Nos. 3,011,920 and 3,532,518, requiring an acceleration step, and the one-step catalyst compositions of U.S. Pat. No. 4,863,758, which do not require acceleration). This same method is used to activate tungsten circuitry surfaces for subsequent electroless nickel plating.\nFor copper surfaces, it is also known to activate them by application of a prestrike layer of nickel from an electroless bath utilizing a boron compound, e.g., dimethyl amine borane or borohydride, as the reducing agent, which prestrike layer then serves to catalyze electroless nickel deposition from a more conventional bath (e.g., hypophosphite-reduced).\nThe palladium-based activators for electroless nickel plating of copper or fused tungsten surfaces are relatively expensive and must be maintained at carefully controlled concentrations in order to be effective. More problematic is the fact that for printed circuits having copper surfaces which are to be activated for electroless plating in this manner, the procedure involves immersion of the entire circuit board, including insulating areas, into the palladium-containing sol or solution. Although post-rinsing procedures are employed, the catalytic species may nevertheless become entrapped into insulating areas, and will there catalyze electroless nickel plating when the board is immersed in an electroless nickel bath. The resultant unwanted areas or paths of metal on the insulating areas can lead to undesirable cross-talk or shorting between conductive circuitry on the board. This same problem occurs with tungsten-ceramic packages, i.e., the palladium-based activator can become entrapped or adhered to insulating ceramic surfaces and promote undesired metallization thereof.\nThe same problem occurs with the prestrike method of activation for copper surfaces, i.e., the very activity of the borane or borohydride-reduced electroless nickel bath which enables it to deposit a strike layer of nickel over the copper also can result in undesired deposition of nickel on insulating areas of the board. Also in common with the palladium-based activators, these nickel strike baths are relatively expensive and their composition must be carefully controlled to avoid spontaneous plating of all metal in the solution."}
{"text": "An HVDC power transmission network (also called an HVDC power distribution system or meshed HVDC grid) uses direct current (DC) for the transmission of electrical power, in contrast to the more common alternating current (AC) systems. For long-distance transmission, HVDC systems may be less expensive and may suffer lower electrical losses. In general, an HVDC power transmission system comprises at least one long-distance HVDC link or cable for carrying direct current a long distance, e.g. under sea, and converter stations for converting alternating current to direct current for input to the HVDC power transmission system and converter stations for converting direct current back to alternating current.\nIn an meshed or, highly meshed HVDC grid, the DC node voltage might not be sufficient to control load flow congestion. A series connected converter may alleviate this problem by adding a fictitious resistance in the line. If the injected series voltage is VAB, to ensure that the cable current is Idc, a fictitious series resistance Rinj=VAB/Idc is injected. Insertion of the series converter in a congested HVDC network is shown in FIG. 1 where a series converter is connected between A and B.\nIn a voltage source converter (VSC) based DC grid, the current direction changes based on the direction of power flow. Hence, for complete control of load flow, the series converter works in all the four quadrants as illustrated in FIG. 2.\nThe AC-DC converters in FIG. 1 represent DC Grid converter stations typically with asymmetrical monopoles with separate converter for positive and negative polarity, or they can be balanced bipolar converters. Both of the converter topologies are shown in FIG. 3. In FIG. 3, for both of the converter types the series converter location is shown. For symmetric voltage profile in the positive and negative lines a series converter (shown as SC in FIG. 3) is positioned in both positive and negative lines as shown in FIG. 3. The aim is to change the load sharing between the two cables x and y (both are positive lines) or, between the two cables z and q (both are negative lines) which is achieved by the series converter.\nThe series converter can be a voltage source converter (VSC). However, the harmonic content of the series injected voltage must be within permissible limits. In case the series converter is powered from an AC source (local AC grid, transformer tertiary winding) the power quality also has to be maintained at the AC side, such that the AC source is not polluted due to the series converter. These conditions may necessitate filters both in the DC and the AC side of the series converter.\nSince the converter is connected in series with the HVDC line, rated current can flow through it, although the voltage produced by the converter is relatively small. Multi-level voltage source converters can reduce the voltage across each of the switching device which forms the converter, but rated line current still flows through them.\nThe protection of the converter is also an important criterion. VSCs are naturally vulnerable to DC side faults (faults on the HVDC cable). This necessitates requirement of fast and costly protection equipment so as to bypass the series converter efficiently at the event of a fault. Additional reactors might also be required in the main HVDC line to limit cable discharge current through the series converter at the time of the fault (pole-ground/pole-pole). These reactors need to be designed for rated current, thereby significantly adding to the cost of the overall solution.\nSeries converters of different known designs are described in WO 2012/037966, WO 2012/037964 and DE 1513827."}
{"text": "Many computer applications such as banking and ticket reservation systems require a fault-tolerant computer that is continually on-line. Thus, the connection and disconnection of modules to a bus must be done without shutting down the system.\nWhen a load is connected to an energized bus a voltage drop caused by the sudden appearance or disappearance of the load can cause damage to data or other components of the system. Further, when the load is disconnected from an energized bus arcing across the contacts damages and reduces the useful life of the contacts. Additionally, a fault in one component of the system can cause excess current flow that may damage the component.\nAccordingly, a need exists for obviating the problems caused by connecting and disconnecting loads to an energized bus for preventing faults from affecting an entire system."}
{"text": "A mechanic servicing or repairing an automobile often utilizes a shop manual or other technical books while working and must refer to the book repeatedly at the various stages of the repair process. While working on the engine, for example, from above, it is possible to support the manual on a fender or a workbench, which is not always an ideal arrangement, but when working beneath the car he is generally on his back, and to refer to a manual laid on the floor from this position, often with his hands occupied, is extremely awkward if not impossible. The mechanic is thus in need of a mobile book stand which will support a shop manual or the like on an incline alternatively at about shoulder level when he works on the car from above, and at, or slightly above, floor level when he is working beneath the car."}
{"text": "This invention relates to rubber compositions having vibration damping and noise insulating properties, and being suitable for use as engine mounts.\nThe function of engine mounts is to insulate the transmission of vibration from the engine to the frame and passenger compartment of an automobile. The operating temperature of the engine mount can vary in the range from about 25.degree. C. to about 80.degree. C. Therefore, materials having relatively high damping coefficients in this temperature range are desired for engine mount applications. The materials used for engine mounts should also have, preferably, noise insulating properties. Engine vibration is particularly troublesome in the low frequency region (10-20 Hz) while the high frequency region (above 75 Hz) contributes to the generation of noise in the passenger compartment.\nlt is desirable that the damping coefficient of the composition (also referred to as tan delta), which is indicative of the vibration damping character of the composition, be relatively high (i.e. highly absorbing) and uniform at least in the 25.degree.-80.degree. C. operating range, and secondly that the ratio of dynamic to static moduli, which is indicative of the noise transmitting property, should remain relatively low. It is to achieving this balance of properties that this invention is directed.\nVarious rubber compositions have been proposed for engine mount applications. Most commonly, blends of natural rubber with halogenated rubber, e.g. brominated butyl rubber, have been used. Natural rubber is a good noise transmission insulator but a poor vibration damper. Bromobutyl rubber is, on the other hand, a poor noise insulator while it has good vibration damping qualities. The above blends are, therefore, a compromise between these two properties.\nVibration damping is usually considered to be the primary function of engine mounts and noise insulation the secondary function."}
{"text": "The present invention is directed to a patch panel and strain relief bar assembly and, more particularly, to quick release brackets for removably securing a strain relief bar to a patch panel.\nU.S. Pat. No. 6,818,834 discloses two mounting elements used to secure a strain relief bar to the rear side of a patch panel with threaded screws. Moreover, U.S. Pat. No. 6,568,542 discloses two anchor elements used to secure a strain relief bar to the rear side of a patch panel. However, neither of these prior art patents disclose quick release brackets for removably securing a strain relief bar assembly to a patch panel."}
{"text": "The present invention relates to a tube clamping structure in which a tube and a tubular opposing member such as a rigid pipe made, for instance, of resin, metal, etc. are connected in a fluid-tight condition. Such a structure is favorably used for clamping structures between tubes and metallic pipes for use in automobiles calling for resistance to gasoline.\nA conventionally known clamping structure of this type is disclosed in Japanese Patent Unexamined Publication No. 4-171391 (1992). More particularly, this is a clamping structure as illustrated in FIG. 10 between a thin wall tube 7 and a metallic tube 5 arranged in the following manner: a connecting portion of the thin wall tube 7 is fitted to an outer periphery portion of a connecting portion of the metallic tube 5; a protecting sleeve 6, which is made of the same material as the thin wall tube 7, is fitted to an outer periphery portion of the connecting portion of the thin wall tube 7; and the outer periphery portion of the connecting portion of the protecting sleeve 6 is swaged by a metallic sleeve 8 at least at two spots that are located remote from each other. The structure is further arranged in that at least one groove 5a is provided on the outer periphery portion of the metallic tube 5, the groove 5a being located on either one of the at least two swaging portions and assuming a dimension such that its depth dimension is smaller than half a thickness of the thin wall tube and its width dimension is smaller than a width dimension of the swaging portion located at one spot of the metallic sleeve 8.\nHowever, since relatively hard resin was employed for the thin wall tube 7, drawbacks were presented in that sealing properties between an inner peripheral surface of the tube and an outer peripheral surface of the metallic tubes 5 were inferior. While structures are accordingly employed for strongly swaging the sleeve from the outer periphery in order to improve sealing properties, it is difficult to employ such structures in assembly worksites for automobiles or the like, since such swaging operations need to be performed by using a swaging device. Problems of leakage of fluid were also presented that were due to creases generated in the course of swaging.\nThe present invention has been made in view of these circumstances, and it is a first object of the present invention to provide a tube connecting structure capable of easily improving sealing properties between tubes and tubular opposing members in addition to enabling easy connection at any spot, e.g. assembly worksites for automobiles. It is another object of the present invention to provide a tube connecting structure for preventing leakage of fluid owing to creases generated in the course of swaging.\nFor achieving the above objects, the present invention provides a tube clamping structure for connecting an opposing member to one end of a tube comprising:\na tube extending from one end to another end; an engaging end portion being formed at one end of the tube and extending in an axial direction, and including an annular concave portion that is formed at an inner peripheral surface thereof and that is located to be remote from an one end surface in the axial direction; a tubular elastic sealing member inserted and attached to the engaging end potion of the tube; a tubular opposing member comprised with an inserting end portion that can be inserted into an inner periphery of the elastic sealing member and including an annular protrusion projecting outward from its outer peripheral surface in a radial direction and being located remote from a tip end thereof; and a clamping member annularly disposed on the outer periphery of the tube, wherein the elastic sealing member is comprised of an inner tube portion with an outer peripheral surface that is fitted to an inner peripheral surface of the engaging end portion of the tube and an outer tube portion provided at one end side of the inner tube portion, being bent outwardly over as to cover the one end surface of the engaging end portion while covering the outer peripheral surface of the engaging end portion, wherein the outer peripheral surface of the inner tube portion is comprised with an annular fitting convex portion matching the annular concave portion of the tube, and the inner peripheral surface thereof with an annular fitting concave portion for fitting the annular protrusion of the opposing member, and wherein the clamping member is clamped on the outer periphery of the outer tube portion with the inserting end portion of the opposing member being inserted and attached to the inner peripheral surface of the tubular elastic sealing member.\nIn the present invention, the opposing member is held against the tubular elastic sealing member so as to be immovable in the axial direction owing to the above-described arrangement in which the tube is comprised with the engaging end portion with an annular concave portion being formed at a position being remote from one end surface thereof in the axial direction, in which the tubular elastic sealing member is comprised of an inner tube portion with an outer peripheral surface and fitted to an inner peripheral surface of the engaging end portion of the tube, and an outer tube portion provided at one end side of the inner tube portion, being bent outwardly over as to cover one end surface of the engaging end portion while covering the outer peripheral surface of the engaging end portion, and in which the outer peripheral surface of the inner tube portion is comprised with an annular protrusion matching the annular concave portion of the tube and the inner peripheral surface thereof with a fitting concave portion for fitting the annular protrusion of the opposing member. The tubular elastic sealing member is further held against the tube so as to be immovable in the axial direction.\nOne end of the tube is pinched and held between the inner tube portion and the outer tube portion of the elastic sealing member and is clamped by a clamping member from the outer periphery of the outer tube portion. Thus, it is possible to prevent the opposing member from slipping off the tube.\nIt is also possible to provide an annular convex portion at an intermediate portion of the engaging end portion as to project outward from the outer peripheral surface of the engaging end portion in the radial direction. The annular convex portion serving to secure creases that are generated at the engaging end portion in the course of clamping, it is enabled to improve sealing properties between the engaging end portion of the tube, the inner tube portion of the elastic sealing member, and the inserting end portion of the opposing member.\nIt is alternatively possible that employ an arrangement in which the annular concave portion formed on the inner peripheral surface of the engaging end portion is formed at an intermediate portion of the engaging end portion, in which the clamping member is disposed on a location closer to the one end than the annular convex portion, and in which the inner tube portion of the elastic sealing member is provided with an extension extending inward from the annular fitting convex portion of the inner tube portion in the axial direction. By the provision of the extension at the inner tube portion of the elastic sealing member, sealing between both end of the engaging end portion can be reliably performed without being affected through clamping with the extension being located remote from the portion at which clamping through the clamping member is performed.\nA portion of the engaging end portion of the tube ranging from the annular concave portion to the one end surface may be formed to assume a wave-shaped tubular portion with waveforms being formed in a peripheral direction. With this arrangement, it is possible to secure sealing properties at the clamping portion and to reliably prevent leakage of fluid because generation of creases can be smoothed out through the wave-shaped tubular portions at portion where clamping of the clamping member is performed.\nAt least one slit extending from the one end surface in the direction of the annular concave portion may be provided at the engaging end portion of the tube in the peripheral direction thereof. With this arrangement, it is possible to secure sealing properties at the clamping portion and to reliably prevent leakage of fluid because generation of creases can be smoothed out through the slits where clamping of the clamping member is performed. Note that it is easier to form slits rather than wave-shaped tubular portions.\nAn inner end portion of the inner tube portion in the axial direction may be branched at an intermediate position in a thickness direction to comprise an annular outer branch portion and an inner branch portion that respectively extend outward and inward with respect to the radial direction in an inclined manner. With this arrangement, sealing properties can be favorably secured with the inner tube portion closely contacting the engaging end portion at the outer branch portion, while the inner branch portion closely contacts the opposing member for favorably securing sealing properties.\nThe inner tube portion may be alternatively comprised of a base portion extending outward in the axial direction from the annular fitting concave portion and an extension extending inward in the axial direction from the fitting concave portion, wherein the base portion assumes an inner diameter that is slightly larger than an outer diameter of the inserting end portion of the opposing member while the extension assumes an inner diameter that is slightly smaller than the outer diameter of the inserting end portion of the opposing member. With this arrangement, it is possible to readily insert the inserting end portion to the base portion while improving sealing properties between the elastic sealing member and the inserting end portion in the extension.\nIt is also possible to provide annular projecting sealing portions that coaxially project outward and inward in radial directions on a plurality of spots on the outer peripheral surface and the inner peripheral surface of the extension of the inner tube portion in the axial direction. With this arrangement, sealing properties can be favorably secured with the projecting sealing portion projecting outward of the inner tube portion closely contacting the engaging end portion, and further favorably securing sealing properties with the projecting sealing portion projecting inward of the inner tube portion closely contacting the opposing member.\nThe tube may be arranged either as a single-layered structure of resin, a double-layered structure with an inner layer of resin, conductive resin or metallic film and an outer layer of resin, or a triple-layered structure with an adhesive layer being formed intermediate of the inner layer and outer layer. With this arrangement, it is enabled to select tube structures in accordance with various purposes.\nThe elastic sealing member of the present invention may be alternatively arranged to comprise an enveloping portion covering only the one end surface of the engaging end portion of the tube and an inner tube portion including an outer peripheral surface that is fitted to the inner peripheral surface of the engaging end portion, an annular fitting convex portion formed on the outer peripheral surface of the inner tube portion that matches the annular concave portion of the tube, and an annular fitting concave portion formed on the inner peripheral surface of the inner tube portion that fits the annular protrusion of the opposing member, wherein the clamping member is clamped on the outer periphery of the engaging end portion when employing such a structure. By omitting the outer tube portion of the elastic sealing member in this manner, it is possible to readily insert the elastic sealing member to the tube while still exhibiting the above effects."}
{"text": "It is well known in the art of integrated circuits that manufacturing processes are imperfect—on each wafer some, but not all, integrated circuits function fully. It is also known that defects tend to occur in clusters. Therefore, for circuits having multiple, large, functional units, there will be a substantial population of manufactured integrated circuits where one functional unit is defective, or even a small portion of a functional unit is defective, but other functional units on the same integrated circuit function properly.\nThe probability that an integrated circuit will have one or more defects increases as the size of the integrated circuit increases. Further, the cost of fabrication increases as integrated circuit size increases. A high performance single or multiple-processor integrated circuit can be quite large. It is therefore desirable to find ways of selling at least some of those integrated circuits that contain one or a few defects.\nTypically, integrated circuits are tested before they are packaged. Those processor circuits that have defective units are often discarded before packaging; their fabrication cost is wasted but further investment in them is prevented.\nSolutions that permit selling partially-defective integrated circuits are known for memory integrated circuits. For example, memory integrated circuits often have spare blocks of memory that can be substituted for defective sections of the circuit. Before spare blocks became common, memory circuits were occasionally packaged with a high-order address bit wirebonded to a constant, then sold as memories of half capacity, with defective sections of the chip disabled. Wirebonding a bondpad to a constant to configure an integrated circuit is known as a bonding option; bonding options may also be implemented through a package trace where a particular package lead is tied to a particular logic level. Bonding the high-order address bit to power or ground permitted sale of packaged circuits having partial capacity, with the same circuit pinout for defects in either the high or low half of the memory array.\nModern integrated processor circuits of high performance are fabricated with at least some cache memory on the processor integrated circuit. Some of these circuits have been designed with bonding options such that a portion of cache may be disabled; a technique that permits product differentiation as well as sale of partially defective circuits. Some of these circuits also have spare blocks of memory that can be substituted for defective sections of cache.\nIt is known that some available processors have multiples of certain functional units. For example, a processor integrated circuit may have more than one integer execution unit, or more than one floating point execution unit. These units are referenced herein as redundant units. Historically, processors with one defective integer unit, or one defective floating point unit, are typically discarded. Similarly, if multiple processor integrated circuits are sold with a defect in an integer unit or floating point unit of one processor, the entire defective processor is disabled.\nIt is also known that some processors have functional elements that enhance performance but which are not essential to processor operation. For example, deletion of a branch prediction unit, speculative prefetch unit, or a speculative execution unit may degrade but not kill performance. These functional elements are referenced herein as performance-enhancing units.\nBranch prediction and speculative execution are common techniques for minimizing the impact of conditional branch instructions on processor performance.\nBranch prediction involves using dedicated hardware to guess whether conditional branches will be taken or not taken. This may be done by tracking recent results of particular conditional branch instructions, and predicting that a particular branch instruction will cause a branch if it caused a branch when it was last executed.\nSpeculative execution involves a machine's beginning execution of instructions after a conditional branch is encountered in an instruction stream when it is not yet known if the instructions are among those that should be executed. Typically, speculatively executed instruction results are not made permanent until it is determined whether they should have been executed. If instructions are speculatively executed, and it is determined that they should not have been executed, their results are discarded.\nIntegrated circuit fabrication processes are known that are capable of producing read-only memory cells that can be programmed after wafer fabrication is complete. These processes have been used to create programmable memory devices, programmable logic devices, and other circuits. Processes of this type that use fusible links or laser-burnable links are known. Other processes are known that create programmable cells that are programmed by tunneling charge onto a charge-trapping layer, such as a floating polysilicon gate or an oxide-nitride boundary, in the gate region of a device.\nIt would be desirable to recover some revenue by using post-fabrication programmable cells to selectively disable defective functional units on processors that have multiples of functional units of those types, and selling the partially defective devices as lower performance processors.\nIt is also known that different application programs have different processor requirements. For example, some programs have much greater floating point computational requirements than others. Similarly, the percentage of branch instructions may also vary significantly.\nBuilt In Self-Test (BIST) is a design technique wherein circuitry of a functional unit is designed, and some additional circuitry added to the unit, such that the unit can effectively test itself."}
{"text": "Gas turbine engines, such as those used to power modern commercial aircraft or in industrial applications, include a compressor for pressurizing a supply of air, a combustor for burning a hydrocarbon fuel in the presence of the pressurized air, and a turbine for extracting energy from the resultant combustion gases. Generally, the compressor, combustor and turbine are disposed about a central engine axis with the compressor disposed axially upstream of the combustor and the turbine disposed axially downstream of the combustor.\nIn operation of a gas turbine engine, fuel is combusted in the combustor in compressed air from the compressor thereby generating and high-temperature combustion exhaust gases, which pass through the turbine. In the turbine, energy is extracted from the combustion exhaust gases to turn the turbine to drive the compressor and also to produce thrust. The turbine includes a plurality of turbine stages, wherein each stage includes of a stator section formed by a row of stationary vanes followed by a rotor section formed by a row of rotating blades. In each turbine stage, the upstream row of stationary vanes directs the combustion exhaust gases against the downstream row of blades. Thus, the blades of the turbine are exposed to the high temperature exhaust gases.\nEach turbine blade typically has an airfoil-shaped hollow body having a concave surface and a convex surface extending between a leading edge of the blade body to a trailing edge of the blade body. The blade body extends generally radially outwardly from a blade root, whereat the blade root is attached to the turbine rotor disk by a dovetail joint, to a blade tip at the distal end of the blade body. In operation, in order to reduce the passage of combustion exhaust gases outside the blade tips as the blades rotate, thereby reducing turbine efficiency, a tight clearance is established by having the blade tips pass in extremely close proximity to the turbine casing or to actually contact the rub surface of a blade outer air seal. In either case, over time the blade tips of the rotating turbine blades are subject to wear from contact with either the engine casing or the rub surface of the blade outer air seal.\nAs a result of the associated physical wear and also oxidation due to expose to the high-temperature of the combustion exhaust gases, the blade tips erode over time in service, the turbine blades actually become shorter. As a consequence, the tip clearance at cruise becomes larger and turbine efficiency is degraded. Therefore, it is customary to take gas turbine engines out of operation for overhaul as necessary to service various parts of the engine. As part of the servicing of the engine, it is conventional practice to inspect the turbine blade tips and remove blades having excessively eroded blade tips. Because turbine blades are usually made of expensive superalloys in order to withstand the high temperatures to which the blades are exposed, and because turbine are often cooled through internal cooling air passages, the presence of which make the blades very expansive to manufacture, it is customary to restore the blades, rather then simply scrapping the removed blades. The removed blades are restored by rebuilding the eroded blade tip sufficiently to return the blade body to its original design length using various techniques for depositing repair alloys unto the eroded tip of the removed blade, thereby salvaging the blade. However, conventional restoration methods require the damaged blades to be removed from the engine, restored, and then replaced in the engine, which necessarily requires the engine to remain out of service for an extended period."}
{"text": "1. Field of the Invention\nThe present invention relates to a nonvolatile semiconductor memory device including two-transistor type memory cells and its manufacturing method.\n2. Description of the Related Art\nIn a nonvolatile semiconductor memory device, a write operation can be carried out after the device is mounted on a printed circuit board. One typical example of such a device is a flash memory which has an advantage in that data storage is possible without a backup battery, and which is highly integrated. Generally, one memory cell of the flash memory is formed by a silicon substrate, a floating gate over a channel region of the substrate, and a control gate over the floating gate.\nRecently, in order to increase an ON current flowing through the memory cell, two-transistor type memory cells have been developed. A prior art memory cell of this type is formed by one selection transistor and one memory transistor connected in series and controlled by voltages of a pair of word lines, and the thickness of a gate insulating layer of the selection transistor is the same as that of a gate insulating layer of the memory transistor. Also, the threshold voltage of the selection transistor is the same as that of the memory transistor where no charges are injected into a floating gate. This will be explained later in detail.\nGenerally, the thinner the gate insulating layer, the larger an ON current flowing through each of the selection transistor and the memory transistor. In the memory transistor, however, if the gate insulating layer is too thin, electrons accumulated in the floating gate are leaked through the gate insulating layer into the silicon substrate. Therefore, it is impossible to greatly reduce the thickness of the gate insulating layer for the memory transistor. On the other hand, in the selection transistor, since the floating gate is not provided, it is possible to greatly reduce the thickness of the gate insulating layer for the selection transistor.\nTherefore, in the prior art memory cell, it is impossible for the gate insulating layer to be optimum both for the selection transistor and the memory transistor.\nFurther, if the threshold voltage of the memory transistor where no electrons are injected into the floating gate is too large, a read operation performed upon this memory transistor may invite a soft write operation, which is called a read disturb phenomenon. Therefore, in order to avoid this read disturb phenomenon, it is preferable that the threshold voltage of the memory transistor be smaller.\nIn the prior art memory cell, however, it is impossible to adjust the threshold voltage of the memory transistor, independent of the threshold voltage of the selection transistor."}
{"text": "1. Field of the Invention\nThe present invention relates to throttle control devices that have a motor and a reduction gear mechanism that is driven by the motor in order to rotate a throttle valve for controlling a flow rate of intake air supplied to an engine, e.g., an internal combustion engine of an automobile.\n2. Description of the Related Art\nJapanese Laid-Open Patent Publication No. 6-264777 teaches a known throttle control device. As shown in FIG. 14, the known throttle control device has a motor 92 and a reduction gear mechanism 94 that is driven by the motor 92 in order to rotate a throttle valve 96 for controlling a flow rate of intake air. A movable section 102 of a throttle sensor 100 is coaxially mounted on one end of a rotary shaft 92s of the motor 92. The movable section 102 has a disk-like configuration including concave and convex portions. The concave and convex portions are formed on the outer periphery of the movable section 102 and are arranged at predetermined intervals in the circumferential direction. A fixed sensing section 104 of the throttle sensor 100 is mounted on the throttle body 91 and is adapted to detect the concavity or the convexity of the movable section 102.\nTherefore, as the movable section 102 of the throttle sensor 100 rotates together with the rotary shaft 92s of the motor 92, the fixed sensing section 104 of the throttle sensor 100 detects the concave or convex portions of the movable section 102 in order to count the number of concave or convex portions moving past the sensing section, so that the rotational angle of the motor 92 and consequently the degree of opening of the throttle valve 96 can be determined. Because the rotational angle of the throttle valve 96 is determined based upon the rotational angle of the motor 92, the accuracy of the measurement of the rotational angle of the throttle valve 92 can be improved in comparison with an arrangement in which the rotational angle of a throttle valve is directly detected.\nHere, in order to provide a level of precision for the measurement, the outer diameter of the movable section 102 is set to be substantially equal to the outer diameter of the motor 92.\nHowever, the throttle sensor 100 of the known throttle control device is configured to detect the concave or convex portions formed on the outer periphery of the disk-like movable section 102 and to count the number of the concave or convex portions in order to obtain the rotational angle of the throttle valve 92. Therefore, the throttle sensor 100 must have a large size in a diametrical direction to accommodate the number of concave and convex portions required for accuracy. For this reason, the space for accommodating the motor 92 having the throttle sensor 100 must be large in size in a diametrical direction in comparison with a space required for accommodating only the motor 92. Therefore, a problem has been that the throttle body 91 must also have a relatively large size"}
{"text": "The current shortages of gas and oil and the high cost of such fuels has stimulated a demand for boilers and furnaces capable of burning comparatively low cost solid fuels such as wood, whereby such boilers and furnaces can be utilized as booster sources of heat in domestic heating systems and, in some cases, may be used as primary heating units.\nThe present invention seeks to satisfy the above need through the provision of a simplified compact and efficient boiler particularly adapted for use in domestic hot water heating systems and also capable of more general utility. The boiler is rectangular in the interest of space conservation and maximized heating capacity. It possesses an upper horizontal slide-in, slide-out power module to facilitate cleaning and maintenance and this module is a rectangular unit consisting of rows of fire tubes assembled in sealed relationship with a water chamber which surrounds the fire tubes in heat exchange relationship therewith. The water chamber connects directly at the rear of the boiler to cold water return and hot water supply lines. An electrical heating element or elements may also be included on the power module to impart to the boiler a secondary method of heating.\nOther important features of the invention are embodied in a unique fire box having simplified and sturdy means to support a cast iron grate and upwardly diverging refractory side walls formed of replaceable fire bricks. An ash removal drawer is arranged below the grate and a horizontal draft tube extending in the back-to-front direction is provided immediately below the grate. A large front access door for the fire box including a draft check makes the feeding of wood into the fire box an easy operation and a separate ash clean-out door with a manual draft device is also provided on the front of the boiler.\nOther important features of the invention will become apparent during the course of the following detailed description.\nThe following U.S. Pat. Nos. are made of record herein under 37 C.F.R. 1.56, but these patents do not possess the above-enumerated features of the present invention coacting in a manner to be described herein for obtaining the operational results of this invention: 1,967,883, 2,553,302, 3,527,260."}
{"text": "The invention relates to apparatus comprising one or more metal electrodes such as electrochemical cells, sensor elements and the like, and more particularly to extending the shelf life of such apparatus.\nMetal electrodes have proved useful in sensor elements for sensing a diverse range of biologically important molecules eg glucose, and for determining physical properties such as pH. A range of possible configurations and applications involving metal electrodes are discussed in our co-pending applications PCT/AU96/00210, PCT/AU96/00365 and PCT/AU96/00723.\nA desirable attribute of all sensor elements is that they have a long shelf lifexe2x80x94that is, the sensing characteristic of the sensor element does not change significantly between manufacture and use (ie on storage).\nIn an electrochemical sensor element the stability of the electrode is critical to the stability of the sensor as a whole. Typically, when left to stand for long periods of time, electrodes become prone to instability in subsequent use thus limiting the useful shelf life. It is thought that such instability is caused by absorption or reaction of the metallic surface with atmospheric contaminants. It has also been observed that filling time of sensors deteriorates on prolonged storage.\nIt is an object of the present invention to overcome or ameliorate at least some of the above disadvantages in the prior art.\nSurprisingly, the present applicant has found that by coating a metal electrode with a monolayer or multilayer of selected materials, electrode behavior can be significantly stabilised in comparison with uncoated metal electrodes without loss of the desirable sensing characteristics of the electrodes.\nAccording to a first aspect, the invention consists in a metal electrode stabilised by a coating, said coating comprising a sulfur containing moiety in its molecular structure, said coating increasing the temporal stability of the electrode relative to a corresponding uncoated metal electrode without modifying other electrochemical properties of said metal electrode.\nxe2x80x9cComprisingxe2x80x9d as herein used is used in an inclusive sense, that is to say in the sense of but not limited to xe2x80x9cincludingxe2x80x9d or xe2x80x9ccontainingxe2x80x9d. The term is not intended in an exclusive sense (xe2x80x9cconsistingxe2x80x9d of or xe2x80x9ccomposed ofxe2x80x9d).\nPreferably, the sulphur-containing moiety is selected from the group comprising thiol, disulphide and SOx. Most preferably the sulphur-containing moiety is a disulphide. The sulphur-containing moiety may also be incorporated in a cyclic structure.\nAccording to a second aspect, the invention consists in a metal electrode stabilised by a coating according to the first aspect, further comprising a hydrophilic group in its molecular structure.\nPreferably, the hydrophilic group is selected from the group comprising hydroxyl, amine, carboxyl, carbonyl, oligo (ethylene oxide) chain, and zwitterionic species. Most preferably, the hydrophilic group is a zwitterionic species. The most preferred zwitterionic species comprises an amine and a carboxyl group in proximity.\nAccording to a third aspect, the invention consists in a metal electrode stabilised by a coating according to the second aspect, further comprising a spacer between the sulphur-containing moiety and the hydrophilic group.\nPreferably, in the third aspect, the spacer consists of an alkyl group or an aromatic group. It is preferable that methylene or ethylene groups be included in the spacer element.\nAccording to a fourth aspect, the invention consists in a method of preparing a metal electrode stabilised by a coating, comprising the step of contacting a metal electrode with a substance comprising a sulphur-containing moiety in its molecular structure.\nAccording to a fifth aspect, the invention consists in a method of preparing a metal electrode stabilised by a coating, comprising the steps of contacting a metal electrode pith a substance comprising a sulphur-containing moiety and a hydrophilic group in its molecular structure.\nAccording to a sixth aspect, the invention consists in a method of preparing a metal electrode stabilised by a coating comprising the steps of contacting a metal electrode with a substance comprising a sulphur-containing moiety, a hydrophilic group and a spacer between the sulphur-containing moiety and the hydrophilic group in its molecular structure.\nThe preferred substances for use in the methods described in the fourth, fifth and sixth aspects are identical to those substances described in respect of the first, second and third aspects.\nThe invention also consists in a method of sensing an analyte, comprising the step of substituting the electrode in a known sensor device with a metal electrode stabilised by a coating according to the present invention, and sensing an analyte.\nVarious embodiments of the invention will now be described by way of example only.\nIt is known in the prior art that thiols form coatings on metals. Thiols have also been used to tether species such as antibodies onto metal surfaces, for instance those of gold particles, for the purposes of immobilisation etc. One would expect that such coatings would also bind contaminants to the surface.\nAs much electrode chemistry involves interaction at the electrode surface, it is thus surprising that coatings used to bind molecules to the metal surface can be useful in preventing contamination of the electrode surface. It is also surprising that notwithstanding the application of the coating an electrode retains desirable electrochemical properties. The procedure for preparing the metal electrode stabilised by a coating involves contacting a metal electrode with selected sulphur-containing compounds, such as thiols, disulphides and compounds of the formula SOx among others being suitable in the context of the present invention. The coatings also desirably contain a hydrophilic group which includes such species as hydroxyl, amino, carboxyl, carbonyl, oligo (ethylene oxide) chains and zwitterionic species. The latter two compounds indicate that compounds having one or more hydrophilic groups are also suitable groups for use in the present invention\nBetween the sulphur group, which acts to tether the molecule onto the metal surface, and the hydrophilic group, which presents a hydrophilic surface, spacers may be employed.\nCompounds useful in the present invention include, but are not limited to 2-mercaptoethanol, 2-mercaptoethylamine, 3-mercaptopropionic acid, thiophene, 3-carboxythiophene, cysteine, homocysteine and cystine. Most preferably the molecule is cystine. In any of the above aspects, the D or L isomers can be used or mixtures of D and L isomers can be used, where such isomers are possible.\nThe compound in accordance with the invention is then applied as a monolayer or multilayer onto the surface of the electrode. It is possible to apply the compound by simply exposing the electrode to the coating material, with the coating material in either the vapour phase or in solution. The substance can be applied by dipping, spraying, painting, printing etc. After application, it is possible to wash the surface of the contacted electrode.\nIn a further aspect of the current invention the layer of the sulfur containing compound can optionally be overcoated with a surfactant layer. The surfactant layer can be applied after the application of the sulfur containing layer or at the same time as the sulfur containing layer, for example the sulfur containing species and the surfactant can be placed in a coating bath into which the electrode material is immersed. Due to the higher affinity of the sulfur containing species for the electrode material it will bind to the electrode surface in preference to the surfactant, leaving the surfactant in a layer over the sulfur containing layer. An example of a suitable surfactant is Triton X-100."}
{"text": "1. Field of the Invention\nThe present invention relates to information distribution. More particularly, the invention relates to the usage of media with content for the distribution of remotely updateable information.\n2. Description of the Prior Art\nThe distribution of multimedia content and other data has traditionally been achieved by using media, such as the floppy disk, magnetic tape or compact disc (CD), containing information supplied by the content provider. Typically, the media are pre-recorded with information from the content provider and, after being sent to the end user, the end user views the information by reading the media. The information is meant to be static for the end user, that is, after the information is pre-recorded onto the media, the information is not changed or updated by the content provider. Such change or update would be impractical, as it would require the end user to return the media to the content provider for updating.\nFor instance, a mail-order retail content provider might pre-record a sales catalog onto a CD and send the CD to a consumer. The consumer would then access the CD and shop for items in the retailer's catalog. If the retailer were to update the catalog in any way (e.g., error connection, price update, cancelled inventory, etc.) the consumer's CD would be out of date. For write-once media used for content distribution, such as standard CD and DVD, distributed information cannot be updated. The retailer would have to send the consumer a new CD with the new catalog information.\nIn the case of re-writable media (e.g., floppy disks, CD-RW, key-chain RAM, etc.), the update process is possible. The end user would return the re-writable media to the content provider and the content provider would over-write the old information with new information. However, due to economic constraints, information-containing media are typically distributed in large quantities. It would be impractical, if not impossible, for content providers to update their distributed re-writable media in this manner.\nWith the advent of the Internet and computer networking, data can be “downloaded”; that is, transferred from one computer to another computer over the network. The Internet and related browser applications simplified this process by allowing content to be served anywhere, anytime, by using HTTP (Hyper-Text Transport Protocol). In the early development and use of HTTP, most of the content on the Internet was “static” in that it was not based on a constantly updated information source, such as connected to live database. With the increased use of scripting and web-based database access/query technologies such CFML (Cold Fusion Markup Language), ASP (Active Server Pages), Perl (Practical Extraction and Report Language) PHP (PHP: Hypertext Preprocessor), and Java Servlets, the content served from one computer to another over the network can be “live” or “dynamic” in the sense that it can be updated in real-time via a database or other application program on the server. In this way, the end user potentially gets a new or different experience each time they are served content from a particular computer. For instance, a website with the weather forecast might change from hour-to-hour or minute-to-minute. Such information would not be useful if pre-recorded on a fixed media.\nTherefore, what is needed is a method and system for continually distributing updateable content to an end user using media that contains static information, without requiring the content provider to re-write the media or to provide updated media to the end user."}
{"text": "Cathodes can emit electrons by photoemission, thermionic emission, and field emission, or as the result of negative electron affinity. A field-emission cathode (or field emitter) provides electrons when subjected to an electric field of sufficient strength. The electric field is created by applying a suitable voltage between the cathode and an electrode, typically referred to as the anode or gate electrode, situated a short distance away from the cathode.\nVarious techniques have been explored for creating field emitters. Chason, U.S. Pat. No. 5,019,003, fabricates a field emitter by depositing preformed electron-emissive objects on a substrate consisting of dielectric and/or electrically conductive material. The preformed objects, which have sharp edges, can consist entirely of electron-emissive material such as molybdenum or titanium carbide. Alternatively, the preformed objects can consist of electrically insulating cores with thin electron-emissive coatings over the insulating cores. The longest dimension of the objects is approximately 1 .mu.m. A bonding layer is employed to bond the objects to the substrate.\nJaskie et al (\"Jaskie I\"), U.S. Pat. No. 5,141,460, discloses a technique in which diamond is used in fabricating a field emitter. Kane et al (\"Kane I\"), U.S. Pat. No. 5,129,850, discloses a related technique for manufacturing a field emitter that utilizes diamond. The fabrication techniques in Jaskie I and Kane I generally entail implanting carbon into a substrate to create diamond nucleation sites and then growing diamond crystallites at the diamond nucleation sites. The resulting regions of diamond crystallites appear to be electron emissive.\nUse of diamond to provide electrons is desirable for a number of reasons. Depending on how it is produced, diamond can have a low work function. This is advantageous because the electric field needed to emit electrons decreases as the work function decreases. Diamond has a low chemical reactivity. In particular, the gases typically present in a sealed vacuum device such as a CRT have little effect on diamond. Also, changes in temperature affect diamond less than most materials used as electron emitters.\nIn Jaskie I and Kane I, the diamond crystallites are grown by chemical vapor deposition (\"CVD\"). While CVD is economically suitable for depositing many materials, diamond CVD is costly because the diamond CVD growth rate is low and a high CVD temperature is needed. The diamond CVD in Jaskie I and Kane I appears too expensive for low-cost volume production of CRTs in flat-panel televisions.\nJaskie et al (\"Jaskie II\"), U.S. Pat. No. 5,278,475, produces a gated field emitter that utilizes diamond crystallites as electron sources. The diamond crystallites are deposited across the upper surface of a supporting structure consisting of a substrate or a patterned layer of conductive/semiconductive material formed on an electrically insulating substrate. A dielectric layer is deposited over the diamond crystallites. A gate (or control) electrode layer, likewise consisting of conductive/semiconductive material, is deposited on the dielectric layer. Openings are formed through the gate electrode and dielectric layer to expose diamond crystallites at selected areas of the supporting structure.\nKane et al (\"Kane II\"), U.S. Pat. No. 5,252,833, discloses a similar gated field emitter in which diamond crystallites provide electrons. The diamond crystallites in Kane II are situated on conductive/semiconductive paths at the bottoms of openings through a dielectric layer and an overlying gate electrode. The diamond crystallites consist of polycrystalline diamond. Taking note of the fact that the (positive) affinity of a material to retain electrons increases the surface work function and thus increases the electric field needed for an electron to escape the material, Kane II indicates that polycrystalline diamond with a (111) crystallographic orientation is particularly useful as an electron source because (111) polycrystalline diamond has a negative electron affinity.\nElectron affinity is an important consideration in choosing an electron source. However, maintaining a negative electron affinity during volume field-emitter production requires special steps. Also, it is not clear that the diamond crystallites in Jaskie II and Kane II will be securely fixed to the underlying material in a manner that permits a control voltage to be suitably impressed on the diamond crystallites. As a result, the gated field emitters of Jaskie II and Kane II may not perform well. It would be advantageous to have an electron-emitting device in which diamond or a related carbon-containing material can be utilized as an electron source and which can be fabricated in a manner that avoids the above-mentioned disadvantages of the prior art."}
{"text": "The subject matter disclosed herein relates to fuel injectors, and, more particularly, to fuel injectors for gasifiers.\nAn integrated gasification combined cycle (IGCC) power plant includes a gasifier with one or more fuel injectors. The fuel injectors supply a fuel, such as an organic feedstock, into the gasifier along with oxygen and steam to generate a syngas. However, existing fuel injectors may have several passages for the fuel, oxygen, and other materials injected into the combustion chamber, thereby increasing the complexity and cost of the fuel injectors. In addition, existing fuel injectors may have narrow operating windows and thus, may be replaced when operating conditions for the gasification system change."}
{"text": "Earth-boring tools for forming wellbores in subterranean earth formations generally include a plurality of cutting elements secured to a body. For example, fixed-cutter earth-boring rotary drill bits (also referred to as “drag bits”) include a plurality of cutting elements fixedly attached to a bit body of the fixed-cutter drill bit. Similarly, roller cone earth-boring rotary drill bits include cones that are mounted on bearing pins extending from legs of a bit body such that each cone is capable of rotating about the bearing pin on which it is mounted. A plurality of cutting elements may be mounted to each cone of such a roller cone drill bit.\nThe cutting elements used in fixed-cutter, roller cone, and other earth-boring tools often include polycrystalline compact cutting elements, e.g., polycrystalline diamond compact (“PDC”) cutting elements. The polycrystalline compact cutting elements include cutting faces of a polycrystalline compact of a polycrystalline material such as diamond or another super hard material (collectively referred to herein as “super hard material”).\nPolycrystalline compact cutting elements may be formed by sintering and bonding together grains or crystals of super hard material in the presence of a metal solvent catalyst. (The terms “grain” and “crystal” are used synonymously and interchangeably herein.) The super hard material grains are sintered and bonded under high temperature and high pressure conditions (referred to herein as “high pressure, high temperature processes” (“HPHT processes”) or “high temperature, high pressure processes” (“HTHP processes”)). The HPHT process forms direct, inter-granular bonds between the grains of super hard material, and the inter-granularly bonded grains form a “table” of the polycrystalline material (e.g., diamond or alternative super hard material). The table may be formed on or later joined to a cutting element supporting substrate."}
{"text": "The present invention generally relates to color video signal recording and reproducing apparatuses, and more particularly to a color video signal recording and reproducing apparatus in which a time base compressed line sequential color difference signal which is obtained by time base compressing a line-sequential color difference signal is time-division-multiplexed with a low-frequency component of a luminance signal so as to obtain a time-division-multiplexed video signal, and a frequency modulated signal which is obtained by frequency-modulating the time-division-multiplexed signal is recorded on a recording medium together with a high-frequency component of the luminance signal, and further, a signal which is reproduced from the recording medium is subjected to a signal processing complementary to the signal processing carried out at the time of the recording so as to obtain a reproduced color video signal which is in conformance with a standard system.\nAmong the existing color video signal recording and reproducing apparatuses such as a video tape recorder (VTR), the popular recording and reproducing apparatus employ the following recording and reproducing method. That is, a luminance signal and a carrier chrominance signal are separated from a composite color video signal of a standard system such as the NTSC, PAL, and SECAM systems. The separated luminance signal is frequency-modulated, and the separated carrier chrominance signal is frequency-converted into a low-frequency band. The frequency converted carrier chrominance signal is frequency-division-multiplexed with the frequency modulated luminance signal, and this frequency-division-multiplexed signal is recorded on a recording medium. At the time of a reproduction, a signal which is reproduced from the recording medium is subjected to a signal processing which is complementary to the signal processing carried out at the time of the recording, so as to obtain a reproduced composite color video signal which is in conformance with the original standard system. In other words, the popular recording and/or reproducing apparatuses employ the so-called low-band conversion recording and reproducing system.\nThe recording and reproducing apparatus employing the low-band conversion recording and reproducing system, has the following advantages. (I) Because the frequency band of the luminance signal can be selected arbitrarily, the apparatus is especially advantageous when applied to a home use VTR in which the frequency band in which the recording and reproduction can be carried out is relatively narrow. (II) A demodulated chrominance signal is uneasily affected by a time base deviation upon reproduction in the VTR. (III) There is little beat interference because only the luminance signal passes through frequency modulating and frequency demodulating systems, and a pilot signal is not recorded and reproduced. (IV) The frequency modulated luminance signal has a high-frequency biasing effect, and enables recording of the carrier chrominance signal with a satisfactory linearity.\nOn the other hand, there is a trend to lengthen the recording and reproducing times of the recording and reproducing apparatus, and there is thus a trend to lower the traveling speed of a magnetic tape which is used as the recording medium. Further, there is a demand to enable reproduction of an audio signal with a higher quality. However, since the audio signal is recorded on and reproduced from the magnetic tape by a stationary head, the relative linear velocity between the magnetic tape and the stationary head is small. For this reason, when the traveling speed of the magnetic tape is reduced, the frequency characteristic of the audio signal becomes greatly deteriorated compared to the frequency characteristic of the video signal which is recorded on and reproduced from the magnetic tape by a plurality of rotary video heads. Therefore, it was impossible to record and reproduce the audio signal with a high quality.\nThus, another recording and reproducing system was proposed. According to this proposed recording and reproducing system, the audio signal is converted into a predetermined signal format, and the audio signal having this converted signal format is superimposed on the video signal. The superimposed signal is recorded on and reproduced from the magnetic tape by the rotary video heads. In this proposed recording and reproducing system, the audio signal is recorded on and reproduced from the magnetic tape by the rotary video heads which have a high relative linear velocity with respect to the magnetic tape. For this reason, even in a case where the traveling speed of the magnetic tape is slow, the audio signal can be recorded and reproduced with a quality which is far superior to the quality with which the audio signal can be recorded and reproduced by the stationary head in a state where the traveling speed of the magnetic tape is not reduced.\nHowever, in the recording and reproducing apparatus employing the low-band conversion system described before, two signals, that is, the luminance signal and the carrier chrominance signal must be transmitted in the narrow frequency band. Hence, the frequency band in which the frequency-division-multiplexed signal made up of the frequency modulated luminance signal and the frequency converted carrier chrominance signal is recorded and reproduced, had to be set to a narrow frequency band. Accordingly, the frequency band in which the luminance signal was recorded and reproduced was narrow, and thus, it was impossible to improve the resolution. Moreover, in the proposed recording and reproducing system in which the audio signal is recorded on and reproduced from the same track as the video signal by the rotary video heads with the low-band recording and reproducing system described before, the recording and reproduction were carried out by arranging a frequency modulated audio signal which is obtained by frequency-modulating a carrier by the audio signal into an unoccupied frequency band between the frequency converted carrier chrominance signal and the frequency modulated luminance signal. Therefore, in this proposed recording and reproducing system, the frequency band of the frequency modulated luminance signal became further limited due to the frequency modulated audio signal, and there was a disadvantage in that the resolution became poorer."}
{"text": "1. Technical Field\nThe present invention relates to a thermal detector, to a thermal detector device, to an electronic instrument, and to a method of manufacturing a thermal detector device.\n2. Related Art\nIn the field of thermal detectors, elements such as thermopiles, pyroelectric elements, and bolometers are known. A thermopile uses thermocouples to directly detect a rise in temperature in a light absorbing film in association with absorption of light.\nA pyroelectric element utilizes the pyroelectric effect of a ferromagnetic body to detect a temperature rise in a light absorbing film in association with absorption of light. For example, ferroelectric PZT (lead zirconate titanate), lithium tantalate, and other such crystals with a high dielectric constant give rise to changes in the level of electrical polarization when heated or cooled. Specifically, spontaneous changes in the level of polarization arise when the temperature changes, producing changes in the level of surface charge; whereas in the absence of temperature change, the surface charge is neutralized. In association with changes in polarization conditions, pyroelectric current flow is produced due to changes in the amount of surface charge between electrodes connected to both ends of the ferroelectric crystal. By detecting this pyroelectric current, the quantity of irradiating light (infrared or the like) can be sensed.\nA bolometer detects temperature rise associated with light absorption, in the form of changes in resistance of a thermosensitive resistance element, for example.\nTypically, a thermal detector has a structure that lacks a cooling system. Consequently, it is necessary to have a structure whereby the element is housed in a hermetic package or otherwise situated in a reduced pressure environment, while providing thermal separation from the substrate and neighboring film to prevent as much as possible the diffusion of heat produced by received light (infrared or the like) to the surrounding area. One effective way to prevent dissipation of heat to the substrate and avoid diminished detection characteristics of a thermal detector is to employ a structure in which, for example, a cavity for thermal separation is provided between the substrate and the thermal detector element (see Japanese Laid-Open Patent Application 2000-205944 and Japanese Laid-Open Patent Application 2002-214038 for example). Japanese Laid-Open Patent Application 2000-205944 discloses a thermal type infrared array sensor having a cavity for thermal separation, and Japanese Laid-Open Patent Application 2002-214038 discloses a pyroelectric type infrared detector element having a cavity for thermal separation.\nAdditionally, infrared detector elements, which are one type of thermal detector element, are employed in the field of small scale elements as personal sensors for example, and in the field of large scale arrays as infrared camera devices for example. While initially developed as military technologies, in recent years application in consumer products is progressing, and various applications for infrared detection may be anticipated in the future."}
{"text": "The present invention relates to nanopore/nanochannel devices, and more specifically, to reducing the entropy of molecules via agarose gel in nanopore/nanochannel devices.\nNanopore sequencing is a method for determining the order in which nucleotides occur on a strand of deoxyribonucleic acid (DNA). A nanopore (also referred to a pore, nanochannel, hole, etc.) can be a small hole in the order of several nanometers in internal diameter. The theory behind nanopore sequencing is about what occurs when the nanopore is submerged in a conducting fluid and an electric potential (voltage) is applied across the nanopore. Under these conditions, a slight electric current due to conduction of ions through the nanopore can be measured, and the amount of current is very sensitive to the size and shape of the nanopore. If single bases or strands of DNA pass (or part of the DNA molecule passes) through the nanopore, this can create a change in the magnitude of the current through the nanopore. Other electrical or optical sensors can also be positioned around the nanopore so that DNA bases can be differentiated while the DNA passes through the nanopore.\nThe DNA can be driven through the nanopore by using various methods, so that the DNA might eventually pass through the nanopore. The scale of the nanopore can have the effect that the DNA may be forced through the hole as a long string, one base at a time, like thread through the eye of a needle. Recently, there has been growing interest in applying nanopores as sensors for rapid analysis of biomolecules such as deoxyribonucleic acid (DNA), ribonucleic acid (RNA), protein, etc. Special emphasis has been given to applications of nanopores for DNA sequencing, as this technology holds the promise to reduce the cost of sequencing below $1000/human genome."}
{"text": "1. Field of the Invention\nThe present invention relates to a substrate with a built-in component, an electronic component being built in the substrate with a built-in component, and a method for manufacturing the same.\n2. Description of the Related Art\nA substrate described in, for example, Japanese Unexamined Patent Application Publication No. 2007-305674 is a known substrate with a built-in component. In the following, the substrate with a built-in component described in Japanese Unexamined Patent Application Publication No. 2007-305674 will be described with reference to FIG. 8. FIG. 8 is a cross-sectional view showing the substrate with a built-in component.\nA substrate with a built-in component 101 is formed by stacking a plurality of insulating layers 102, which are made of a thermoplastic resin, and by performing thermocompression bonding thereon. Wiring conductors 103, which are made of copper foil or the like, are formed on some of the insulating layers 102. The wiring conductors 103 are electrically connected by vias 104 to lands 105, which are formed on a main surface of the substrate with a built-in component and on the surface opposite the main surface. An electronic component 106 is built in the inside of the substrate with a built-in component 101. Electrodes 106a of the electronic component 106 are electrically connected to the wiring conductors 103 and the lands 105 by the vias 104.\nIn the case where a plurality of resin sheets composing the substrate are bonded by thermocompression, there is a possibility that wiring conductors and vias are displaced and open defects and short defects occur in such a substrate with a built-in component as a result of the resin flow of a resin sheet in the plane direction or in the thickness direction. A problem has been newly discovered in which the above-described problem occurs more often, particularly in the case where wiring in the periphery of a built-in component is improved to be finer and denser. The resin flow is particularly large in the periphery of the built-in component because the resin in the periphery of the built-in component moves so as to fill the gap between the component and the resin."}
{"text": "1. Field of the Invention\nThe present disclosure relates to a liquid crystal display (LCD) device and a driving method thereof, in which the size of a printed circuit board (PCB) with a driving circuit mounted thereon is reduced and the manufacturing cost is saved.\n2. Discussion of the Related Art\nLCD devices have numerous advantages, for example, advanced manufacturing technology, good drivability of a driving means, low power consumption, high-quality images, and a large screen. Therefore, LCD devices are popular. Also, LCD devices are being applied to various fields such as portable computers including notebook computers, office automation equipment, portable multimedia equipment, indoor/outdoor display devices, etc., and the application fields of LCD devices are continuously expanding.\nLCD devices adjust the light transmittances of respective pixels according to an input video signal, thereby displaying an image.\nFIG. 1 is a diagram illustrating a related art LCD device. FIG. 2 is a diagram illustrating a connection structure between a gamma block and data driver integrated circuit (IC) of the related art.\nReferring to FIGS. 1 and 2, the related art LCD device includes a liquid crystal panel 10 that displays an image by using an input image signal, a backlight unit (not shown) that supplies light to the liquid crystal panel 10, and a driving circuit that drives the liquid crystal panel 10.\nThe liquid crystal panel 10 includes an upper substrate (color filter array substrate), a lower substrate (thin film transistor (TFT) array substrate), and a liquid crystal layer formed between the upper substrate and the lower substrate. The liquid crystal panel 10 includes a plurality of pixels that are arranged in a matrix type, and adjusts the transmittance of light irradiated from the backlight unit to display an image.\nThe driving circuit includes a gate driver (not shown), a data driver (not shown), a gamma voltage generator 40, a timing controller 50, and a power supply (not shown).\nHere, a plurality of data driver ICs 20, the gamma voltage generator 40, and the timing controller 50 are mounted on a PCB 30.\nThe gate driver includes a plurality of gate driver ICs, and sequentially supplies a scan signal to a plurality of gate lines formed in the liquid crystal panel 10 to switch on the plurality of pixels.\nThe data driver includes the data driver ICs 20, and respectively supplies data voltages to a plurality of data lines formed in the liquid crystal panel 10.\nHere, the data driver ICs 20 convert digital image data, supplied from the timing controller 50, into analog data voltages and supply the analog data voltages to the data lines, respectively.\nThe timing controller 50 aligns digital image data inputted from the outside and supplies the aligned data to the data driver ICs 20.\nMoreover, the timing controller 50 generates a plurality of control signals for controlling the gate driver and the data driver, and supplies the control signals to the gate driver and the data driver, respectively.\nAs illustrated in FIG. 2, the gamma voltage generator 40 includes a plurality of gamma blocks 42, and respectively generates a plurality of gamma voltages to the data driver ICs 20.\nA capacitor (not shown), which buffers a gamma voltage and outputs a certain voltage value, is disposed in an output terminal of each of the gamma blocks 42.\nIn FIG. 2, as an example, the gamma voltage generator 40 is illustrated as including ten gamma blocks 42. Each of the gamma blocks 42 includes two resistors that are connected serially between a driving voltage VDD terminal and a ground voltage GND terminal.\nThe gamma blocks 42 generate a first gamma voltage GMA1 to a tenth gamma voltage GMA10 (which have different values) by using two corresponding resistors that are connected serially between the driving voltage VDD terminal and the ground voltage GND terminal, respectively. Furthermore, the gamma blocks 42 supply the first gamma voltage GMA1 to the tenth gamma voltage GMA10 to the data driver ICs 20, respectively.\nHere, the data driver ICs 20 convert the digital image data outputted from the timing controller 50 into the analog data voltages by using positive and negative gamma voltages GMA1 to GMA10 supplied from the gamma voltage generator 40.\nA plurality of transmission lines 60 are formed on the PCB 30, and the gamma voltage generator 40 and the data driver ICs 20 are connected in parallel through the transmission lines 60. The gamma voltages GMA1 to GMA10 generated by the gamma voltage generator 40 are supplied to the data driver ICs 20 through the transmission lines 60.\nIn the related art LCD device having the above-described configuration, when the gamma voltage generator 40 is configured with ten gamma blocks 42, the plurality of transmission lines 60 are required to be formed on the PCB 30 for parallelly connecting the ten gamma blocks 42 and the data driver ICs 20.\nSince the transmission lines 60 are formed on the PCB 30, the area of the PCB 30 increases. Due to this reason, much research is being recently conducted for reducing the area of the PCB 30 on which the driving circuit of the LCD device is mounted. However, since the transmission lines 60 are formed on the PCB 30, there is a limitation in decreasing the area of the PCB 30.\nMoreover, a method of reducing the layer of a PCB is proposed for saving the manufacturing cost of LCD devices, but since the transmission lines 60 are formed on the PCB 30, there is a limitation in decreasing the layer of the PCB 30.\nMoreover, twenty resistors R1 to R20 are required for generating the first gamma voltage GMA1 to the tenth gamma voltage GMA10 by using the ten gamma blocks 42, and ten capacitors are disposed in respective output terminals of the gamma blocks 42, causing the increase in the manufacturing cost of LCD devices."}
{"text": "The present invention relates to mobile communications and, more specifically, to ambient sound-based call functions for mobile communications devices.\nOne of the main conveniences of mobile communications devices is the ability to stay connected while on the go. Many individuals greatly rely on the ability to be in communication with others at virtually any time of the day and at any location. Oftentimes, however, individuals find themselves in physical locations in which noise levels are exceptionally high. When an incoming call is received during this type of situation, the individual has the option to either ignore the call or attempt to conduct the communication over the noise levels. Competing with such a distraction can result in confusion, misinformation, and frustration for both parties to the communication."}
{"text": "A conventional domestic fan typically includes a set of blades or vanes mounted for rotation about an axis, and drive apparatus for rotating the set of blades to generate an air flow. The movement and circulation of the air flow creates a ‘wind chill’ or breeze and, as a result, the user experiences a cooling effect as heat is dissipated through convection and evaporation.\nSuch fans are available in a variety of sizes and shapes. For example, a ceiling fan can be at least 1 m in diameter, and is usually mounted in a suspended manner from the ceiling to provide a downward flow of air to cool a room. On the other hand, desk fans are often around 30 cm in diameter, and are usually free standing and portable. Other types of fan can be attached to the floor or mounted on a wall. Fans such as that disclosed in U.S. D 103,476 and U.S. Pat. No. 1,767,060 are suitable for standing on a desk or a table.\nA disadvantage of this type of fan is that the air flow produced by the rotating blades of the fan is generally not uniform. This is due to variations across the blade surface or across the outward facing surface of the fan. The extent of these variations can vary from product to product and even from one individual fan machine to another. These variations result in the generation of an uneven or ‘choppy’ air flow which can be felt as a series of pulses of air and which can be uncomfortable for a user. In addition, this type of fan can be noisy and the noise generated may become intrusive with prolonged use in a domestic environment. A further disadvantage is that the cooling effect created by the fan diminishes with distance from the user. This means that the fan must be placed in close proximity to the user in order for the user to experience the cooling effect of the fan.\nAn oscillating mechanism may be employed to rotate the outlet from the fan so that the air flow is swept over a wide area of a room. In this way the direction of air flow from the fan can be altered. In addition the drive apparatus may rotate the set of blades at a variety of speeds to optimise the airflow output by the fan. The blade speed adjustment and oscillating mechanism can lead to some improvement in the quality and uniformity of the air flow felt by a user although the characteristic ‘choppy’ air flow remains.\nSome fans, sometimes known as air circulators, generate a cooling flow of air without the use of rotating blades. Fans such as those described in U.S. Pat. No. 2,488,467 and JP 56-167897 have large base body portions including a motor and an impeller for generating an air flow in the base body. The air flow is channeled from the base body to an air discharge slot from which the air flow is projected forward towards a user. The fan of U.S. Pat. No. 2,488,467 emits air flow from a series of concentric slots, whereas the fan of JP 56-167897 channels the air flow to a neck piece leading to a single air discharging slot.\nA fan that attempts to provide cooling air flow through a slot without the use of rotating blades requires an efficient transfer of air flow from the base body to the slot. The air flow is constricted as it is channeled into the slot and this constriction creates pressure in the fan which must be overcome by the air flow generated by the motor and the impeller in order to project the air flow from the slot. Any inefficiencies in the system, for example losses through the fan housing, will reduce the air flow from the fan. The high efficiency requirement restricts the options for the use of motors and other means for creating air flow. This type of fan can be noisy as vibrations generated by the motor and impeller tend to be transmitted and amplified."}
{"text": "This invention relates to area navigation systems, and more particularly to an area navigation system having improved accuracy in the presentation of course deviation information. The invention is described in the particular context of the area navigation systems described and claimed in Assignee's Abnett et al U.S. Pat. No. 3,796,867 entitled, \"ANGLE MODE AREA NAVIGATION COMPUTER\", and in Assignee's Bean U.S. Pat. No. 3,750,942 entitled, \"DISTANCE MODE AREA NAVIGATION COMPUTER\", the disclosures of which patents are incorporated herein to the extent pertinent.\nAn extended description of the concepts of area navigation and its many advantages is set forth in the Abnett et al and Bean patents, and the reader is referred to these patents for more information. Briefly, however, there exists a network of navigation aids called VORTAC's which provide ratio information from which range and bearing information may be derived. In area navigation, the VORTAC range and bearing data are used in conjunction with an arbitrarily selected range and bearing from the VORTAC to a so-called waypoint. From these, an on-board computer generates range and bearing data from the aircraft to the waypoint. Because the location of a waypoint is arbitrary within the range of a particular VORTAC, the pilot is no longer constrained to fly along or near the established airways (defined with respect to a series of VORTAC stations). Instead, one may fly along any convenient course simply by appropriate choice of waypoints. In broad terms, area navigation systems such as those of the Bean and Abnett et al patents facilitate the above by solving a vector triangle defined by the aircraft, the VORTAC, and the waypoint, given the selected waypoint address vector and the measured aircraft to VORTAC vector.\nIn the Abnett et al patent, there is shown an area navigation system in which the aircraft to waypoint vector is used to generate a display of range from the aircraft to the waypoint, and a track angle deviation display, the latter is based on the difference between the actual aircraft to waypoint bearing (the angle component of the aircraft to waypoint vector) and the bearing at which it is desired that the aircraft intercepts the waypoint, i.e. the selected track angle. The angle display is advantageously accomplished by means of a \"left-right\" needle included in the course deviation indicator (or CDI) of the standard VORTAC navigation systems.\nIn the Bean patent, track deviation is displayed on a left-right meter as distance from the selected track.\nEach of the aforementioned types of systems has particular advantages, as described in the respective patents. In both instances, however, one of the most important factors in assuring accurate and useful system operation (apart from good system design) is overcoming a track deviation display instability due to radio propagation defects associated with the VORTAC station, and corresponding errors inherent in the mathematical formulation of the data, and the display mode required for the track error."}
{"text": "When wiring a building, a wire is pulled through a conduit from one location to another. To do this, a small line is pushed or fished from a first side of the conduit to a second side using a wiring system. Common wiring systems may include, for example, a power fish tape, a tugger or puller, a blower/vacuum, a feeding guide pulley, a power feeder, where each of these components accomplishes a separate function.\nThe first side of the conduit and the second side of the conduit are generally disposed within an electrical box. As such, each of the devices associated with the wiring system must be coupled to or supported near the electrical box and generally include their own associated features to be coupled to or supported near the electrical box, making them cumbersome to set up and remove. In addition, electrical boxes vary in size, making some coupling and/or support features for the devices difficult or impossible to use with the various sizes of electrical boxes."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nAn associated service between an Email box and a mobile communication terminal is to bind a user's Email-box number with the user's mobile communication terminal number, so as to make an Email-box user utilize a mobile communication terminal to receive a new-Email reminding message from a bound Email-box, meanwhile the mobile communication terminal may also be utilized to complete various kinds of Email-box operations, such as reading operation, writing operation, replying operation and forwarding operation about an Email.\nFIG. 1 is a schematic diagram illustrating existed exemplary interactive operations between a mobile communication terminal and an Email box. With reference to FIG. 1, the core of interactive operations between the mobile communication terminal and the Email box lies in transit server. The transit server generally needs to complete the following functions.\n1) Receiving a new-Email reminder sent by the Email-box server via the Internet, filtering according to filter-conditions set by user, sending a new-Email reminder which meets the filter-conditions to the mobile communication terminal through Short Message Service (SMS).\n2) The mobile communication terminal initiates requests, such as reading-Email request and/or writing-Email request to the transit server with the SMS mode. After performing a series of checks, such as an authentication check and an integrity check, the transit server forwards the requests to the Email-box server, and provides interactive services, which include a reading operation and a writing operation, for the mobile communication terminal and the Email-box server with the SMS mode.\nHowever, current interactive operation mode between the mobile communication terminal and the Email box has the following disadvantages.\nDue to the single mode for operating Email box, mobile communication terminal users may only complete Email-related operations by accessing an Email box with the SMS mode. The direct result brought about by the single operation mode is inconvenient operations, and low efficiency of interactive operations. Besides, when accessing the Email box with the SMS mode, the operation mode will be loaded down with trivial details, e.g., it is necessary to input instructions frequently, and there may be some limitations for length of the message content, thus efficiency of the interactive operations will be low."}
{"text": "The present invention relates generally to fluidized bed combustion and, more particularly, to certain new and useful improvements in fluidized bed combustion devices and methods for carrying out fluidized bed combustion.\nFluidized bed combustion is currently one of the most promising technologies for advances in the field of combustion furnaces, boilers, incinerators, etc., particularly in these days of exhorbitant energy costs and restrictive environmental protection regulations. A significant reason that fluidized bed combustion holds such promise is my previous invention described in detail in my U.S. Pat. No. 4,240,377 granted Dec. 23, 1980 the disclosure of which is hereby incorporated by reference herein.\nBriefly, fluidized bed combustion permits the use of a wide range of solid, liquid and gaseous fuels, while at the same time reducing pollutants such as nitrous oxide and sulfur dioxide, and providing high combustion and heat exchange efficiencies when compared with conventional combustion techniques. This invention of my '377 patent enables higher combustion and heat exchange efficiencies than previously known fluidized bed combustion devices. Furthermore, it provides several control points to permit immediate variation of steam generating capacity.\nAlthough numerous significant advantages are realized by my invention described in said Pat. No. 4,240,377 (\"my '377 patent\") over the prior state-of-the-art fluidized bed combustion techniques, there are certain aspects which can be improved to provide even further enhanced performance and improved efficiency. For example, the use of the solids circulation system coupled with the deletion of heat exchange tubes in the combustion bed resulted in much higher controlled heat transfer from the bed to the heat exchange surfaces. However, greater efficiency of combustion and heat exchange will enhance such controlled heat transfer.\nIn addition, despite the dramatically low production of sulfur dioxide, there is an economic need to minimize the amount of extraneous sulfur dioxide removing agent (e.g., limestone) added to the combustion chamber while still maximizing removal of the sulfur gas. Furthermore, despite the presence of cyclone gas-solids separators and/or various superheater tubes located above the combustion bed, it is important to retard the loss of fine particulate matter with the flue gas. The best way of doing so is by providing additional means for returning these fine particles to the combustion bed.\nA significant drawback to conventional fluidized bed combustion devices is that limestone and ash must be removed from the combustion bed to maintain an equilibrium of reactivity within the bed. Such removal is usually done by \"purging\" the bed--i.e., by a continuous process of physically removing bed constitutents from the bed for disposal. This results in loss of the heat from the combustion zone as well as unreacted fuel and limestone particles. It also presents a significant disposal problem since dumping can contaminate the aquifer.\nFinally, the efficiency of combustion can be enhanced by increasing the amount of movement, or circulation, of the bed particles. Thus, the greater the circulation of bed constituents within the burning zone, the more efficient the combustion of fuel particles and the greater the reaction to reduce sulfur dioxide.\nAccordingly, it is an object of the present invention to further improve the overall efficiency of fluidized bed combustion and heat exchange in fluidized bed combustion, particularly in a boiler or like apparatus. It is a concommitant object of the invention to provide new and improved fluidized bed combustion devices and methods of carrying out fluidized bed combustion.\nIt is also an object of the present invention to provide a new and improved fluidized bed combustion apparatus and method of carrying out fluidized bed combustion providing heat exchange contact between constituents from the combustion bed and physically separated heat exchange means so as to ensure that the rate of heat exchange may be varied while at the same time constant temperature is maintained within the burning zone.\nIt is a further object of the invention to improve fluidized bed combustion by continuously reducing the size of limestone and ash particles in the combustion bed. As a result, final removal of spent particles in the form of fine dust-like ash and other similarly sized particulate matter collected at the flue gas exit. By the same token, it is also intended that the entrainment of the remaining unspent limestone, fuel particles and other bed constituents with the flue gases be restrained.\nIt is still another object of the invention to improve fluidized bed combustion by generating increased and, therefore, thorough circulation of bed constituents in the combustion bed. At the same time, it is intended to cause even distribution of the bed constituents not only to increase efficiency of combustion but also to enhance reactions in the bed such as the reaction of sulfur dioxide with limestone.\nVarious objects and advantages of the invention have been set forth above and are described in greater detail below. While such improved performance characteristics are specifically described herein, other improvements in various aspects of fluidized bed combustion and heat exchange will be apparent to those skilled in the art either by reason of this description or from practice with the invention as embodied in the novel and improved structures, instrumentalities, methodologies, operational steps and combinations thereof disclosed herein."}
{"text": "Electric routers are the most versatile of all power-driven woodworking machines. A router bit is capable of speeds in the range of from 20,000 to 30,000 rpm and enables the non-professional to turn out work that once was strictly in the province of a master cabinetmaker. In particular, the router is capable of making the most intricate joints, decorative cuts and exacting inlays with a high degree of accuracy. In short, a router is the secret for achieving sophisticated professional touches on wood products that are impossible to duplicate with the use of any other tool.\nIn order to extend the capability of a router. Tables can be purchased to convert the router to a spindle shaper. In particular, the router is normally mounted below the table to have its router bit extend through an opening formed through the top working platen of the table. This arrangement frees the operator's hands for better control of the workpiece and allows operations that would otherwise prove impossible by use of a portable or hand-held router.\nVarious improvements have been made to router tables over the years, including the provision of means for tilting the platen of a router table relative to the rotational axis of the router bit. Routers have also been mounted on a table to swing in an arcuate manner. U.S. Pat. Nos. 3,604,484 and 4,741,370 disclose other types of router tables that have adjustments for the purpose of aiding the operator. In particular, such improvements are primarily intended to increase his or her capacity to perform various work tasks with greater accuracy.\nConventional router tables of the above type are generally expensive to purchase and maintain, complicated to adjust and manipulate and do not always provide the high degree of accuracy required for many sophisticated work tasks."}
{"text": "Photosensitive composition plays an important role in the fabrication of printed circuit boards (PCBs). Photosensitive compositions can be classified, based on their intended applications, into two main categories: photoresists for etching and solder resists. The former are used in the laying out of the copper circuit pattern using photolithographic techniques and will be removed after the copper circuit pattern is completed. The latter are used after the copper circuit pattern is laid out to protect certain portions of the copper circuit from being contacted by the soldering material during the soldering process. Both types of photosensitive compositions can be applied using either a liquid-type or a film-type technique. With the liquid-type technique, the photosensitive composition is applied onto the substrate using screen printing, roller coating, curtain coating, dip coating, etc. With the film-type technique, on the other hand, the photosensitive composition is first coated onto a transparent film, typically a transparent polyester film, which is then overlaid on the substrate and the photosensitive composition is then applied onto the substrate surface by hot-laminating or other suitable techniques.\nPhotosensitive materials, or photoresists in general, can also be categorized according to how they are developed. At the present time, there are two most common types of developments for photoresists: the solvent-development type and the alkaline-development type. The former uses a solvent, typically an organic solvent, to remove (for a negative type photoresist) the exposed portion of the photoresist. On comparison, if the photosensitive composition contains carboxyl groups, the exposed portions can be removed, and thus developed, using an alkaline solution such as sodium carbonate. Because the solvent-type development process can cause solvent exhaust and waste disposal problems which can cause significant increase in the overall manufacturing cost, the current industrial trend is to use the alkaline-development type processes. During the alkaline-development process, the carboxyl-containing polymer molecules will be neutralized by sodium carbonate to form a polymer salt. The unexposed portion of the photosensitive composition, which was not cross-linked, will be emulsified and caused to be dispersed in water. This thus completes the development process utilizing a negative-type photoresist.\nThe photosensitivity of a photosensitive composition is related to the degree of crosslinking after being exposed to UV light. Typically, a greater number of unsaturated double bonds (typically provided by the so-called photomonomers) in the photosensitive composition will result in a higher degree of cross-linking. However, excessive amounts of photomonomers would result in the photosensitive layer being too tacky, thus adversely affecting the image transfer process. With a dry film type photosensitive material, the photosensitive composition is protected by the transparent polyester film from being contaminated and the oxygen-induced polymerization-inhibition, and the requirement for the unsaturated double bonds can be entirely supplied by the photomonomers. Under these conditions, a alkaline-developable photoresist can be prepared by introducing only the carboxyl groups into the polymer binder, there is no need to add double bonds into the polymer binder.\nOn the other hand, when a solution-type photosensitive material is used, because no such protective polyester film is provided, the film thickness and the reactivity of the photosensitive film becomes important, and, as a result, the photosensitive composition can only contain a limited amount of photomonomers. In order to compensate for the decreased photosensitivity due to the reduction in the amount of photomonomers, unsaturated double bonds must be incorporated into the polymer binder so as to increase its post- UV exposure degree of cross-linking. This method has been disclosed in, for example, U.S. Pat. No. 5,055,378, which disclosed a photosensitive composition containing a polymer binder formed from the reaction of an epoxy resin with an carboxyl-containing unsaturated compound so as to impart unsaturated double bonds to the polymer binder. However, the polymer binder disclosed in the '378 patent does not contain carboxyl groups. As a result, it must be developed using a solvent and cannot be developed using an alkaline solution.\nU.S. Pat. Nos. 5,100,767, 5,009,982 and 4,943,516, and Japan Patents 61-243869 and 1-54390 disclosed various photosensitive compositions whose polymer binders are imparted with carboxyl groups by reacting the polymer binder disclosed in the '378 patent with anhydrides. Such an additional step not only allowed the photosensitive composition to become alkaline-developable, it also improved the photosensitivity of the polymer binder by introducing unsaturated double bonds into the polymer binder. However, because the formation of such alkaline-developable photosensitive polymer binder requires excess (about 10 to 20%) amounts of unsaturated carboxyl group containing compounds (such as acrylic acid or methylacrylic acid), which are typically volatile and can cause health concerns, the final finished products can cause health hazards to human bodies, especially after long term exposure to these products. This is one of the main disadvantages of the photosensitive compositions disclosed in the above-mentioned patents.\nU.S. Pat. Nos. 4,764,452 and 5,364,736 and Japanese patents 4-239070 and 2-47657 disclose photosensitive compositions containing polymer binders based on modified styrene-maleic anhydride resins, or the SMA resins. In these photosensitive compositions, partially-esterized SMA resins were reacted with saturated low-molecular weight alcohols in a ring-opening reaction by which carboxyl groups are introduced into the side chains of the SMA resins to form a carboxyl group containing polymer binder which consequently becomes base-developable. Because the reaction to prepare the base-developable polymer binder does not involve unsaturated carboxyl group containing molecules, it does not raise the health hazard problem mentioned above. However, these photosensitive compositions exhibited another type of shortcoming in that, because the polymer binder does not contain unsaturated double bonds, the developed, or hardened, final product does not provide adequate chemical resistance and hardness due to inadequate degree of crosslinking. In order to increase the degree of crosslinking for the post-exposure polymer binder, it has been proposed to add a mono-functional unsaturated alcohol, such as hydroxyethyl (meth)acrylate, in the ring-open reaction of the SMA resin. This method was disclosed in U.S. Pat. Nos. 5,087,552, 5,114,830 and 5,296,334. The use of mono-functional unsaturated alcohols not only imparts the carboxyl group into the SMA resin, it also imparts double bonds so as to increase the degree of cross-linking. However, while this action increases the photosensitivity of the polymer binder, it also returns the old problem of leaving low-molecular weight volatile, smelly and potentially health-hazard molecules in the polymer binder."}
{"text": "This application relates to and claims priority from Japanese Patent Application No. 11-89794 filed on Mar. 30, 1999, the contents of which are hereby incorporated by reference.\n1. Field of the Invention\nThe present invention relates generally to heat exchangers, and is suitably applied to a radiator which radiates heat of coolant of a water-cooled engine of a vehicle into atmospheric air.\n2. Related Art\nJP-U-58-154389 discloses a radiator having plural metal tubes and a header tank communicating with the tubes. The header tank is formed by connecting a metal core plate and a resin tank body by clamping. The tubes are brazed to the core plate. Recently, recycling performance of vehicle parts such as a radiator is demanded to be improved for reducing industrial waste. However, since the above-mentioned radiator is made of at least two kinds of materials including metal and resin, parts of the radiator needs to be divided into metal parts and resin parts for recycling. Therefore, the number of processes for recycling the radiator is increased, and recycling performance of the radiator is low.\nAs shown in FIG. 20, when a header tank 520 of a radiator (not shown) is made of only metal with a relatively small rigidity or a relatively small Young\"\"s modulus such as aluminum, the header tank 520 may not have a sufficient mechanical strength. As a result, the header tank 520 may be readily deformed by an internal pressure thereof to be expanded. Further, a connection portion of the header tank 520 to which a tube 511 is connected has a mechanical strength larger than that of a non-connection portion of the header tank 520 to which no tube 511 is connected. As a result, the non-connection portion of the header tank 520 may be intensively applied with stress and largely deformed.\nFurthermore, as shown in FIG. 21, when the header tank 520 has a rectangular cross-section having a longer side ml and a shorter side m2, a longer-side wall of the header tank 520 has a larger area than a shorter-side wall of the header tank 520. Therefore, the longer-side wall is applied with a larger amount of the internal pressure of the header tank 520 than the shorter-side wall and may be largely deformed. Deformation of the header tank 520 may be restricted by increasing a thickness of a metal plate from which the header tank 520 is formed. However, in this case, a weight and a manufacturing cost of the radiator may be increased.\nIn view of the foregoing problems, it is an object of the present invention to provide a heat exchanger including a metal header tank having a sufficient mechanical strength without increasing a weight and a manufacturing cost of the heat exchanger.\nAccording to the present invention, a heat exchanger has a plurality of metal tubes through which fluid flows, and a metal header tank disposed at a flow-path end of the tubes to extend in a direction perpendicular to a longitudinal direction of the tubes and to communicate with the tubes. The header tank includes first and second walls opposite each other. The first wall has a plurality of connection portions each of which is connected to each of the tubes. The second wall has a plurality of first reinforcement ribs which increase a rigidity of the header tank. The connection portions are arranged in the direction perpendicular to the longitudinal direction of the tubes at a first pitch. The first reinforcement ribs are arranged in the direction perpendicular to the longitudinal direction of the tubes at a second pitch approximately equal to the first pitch.\nTherefore, a rigidity of a non-connection portion of the header tank to which no tube is connected is increased by the first reinforcement ribs, and stress is restricted from being intensively applied to the non-connection portion. As a result, a mechanical strength of the header tank is increased without increasing a thickness of a metal plate from which the header tank is formed, thereby restricting a weight and a manufacturing cost of the heat exchanger from increasing.\nPreferably, the header tank is formed by connecting a first tank member and a second tank member each of which is formed by pressing to have a L-shaped cross-section so that the first and second tank members are formed using the same die."}
{"text": "Traditionally, storage provided in a storage cluster (such as by using a redundant array of independent nodes, or RAIN) is made reliable against hardware failure either through replication of stored objects or erasure coding of stored objects. The former has the advantage that the same unique identifier can access the multiple replicas (using a journal and RAM-based indexing scheme, for example), but has the disadvantage of high bandwidth and storage overhead (depending upon the number of replicas desired, large objects can take up a significant amount of space). The latter enjoys the benefit of a smaller storage footprint and less overhead for similar level of protection against media failures, but suffers from the drawback that each segment of an erasure set is different content that must be separately identified in order to read the object or to reconstruct any lost segments. This identification can be especially problematic when a storage cluster is restarted. Erasure coding will also incur a higher processing overhead and lose its footprint advantage when storing small objects.\nThus, both techniques have disadvantages. Further, some prior art approaches applicable to erasure coding use a control database separate from the storage cluster in order to identify and track segments of a particular object; this approach is problematic because it introduces more overhead and calls into question the availability of this control database and whether or not it needs to be replicated. Also, even though under erasure coding an object can be reconstructed using a subset of the segments used to encode that object (e.g., if there had been a disk failure), it can be time consuming not only to identify which segments are no longer present, but also to locate the remaining segments.\nAccordingly, improved techniques are desired for use with storage clusters in order to take advantage of the benefits of replication and erasure coding as well as to limit exposure after a hardware failure."}
{"text": "The mineralocorticoid hormone aldosterone is produced by the adrenal gland and acts on the distal tubules and collecting ducts of the kidney to increase reabsorption of ions and water in the kidney. Aldosterone causes conservation of sodium, secretion of potassium, increased water retention, and increased blood pressure.\nAldosterone has been implicated in the pathogenesis of cardiovascular diseases such as hypertension and heart failure. In clinical trials, treatment with the nonselective mineralocorticoid receptor antagonist (MRA) spironolactone or the selective MRA eplerenone significantly reduced morbidity and mortality among patients with heart failure or myocardial infarction already taking an angiotensin-converting enzyme inhibitor or a β-blocker. However, significant side effects such as gynecomastia and impotence were observed in male patients receiving spironolactone while hyperkalemia was seen in patients taking either drug."}
{"text": "The present invention relates to information processing techniques, and more particularly to techniques that enable two or more parties to discover mutually held information without disclosing the private or confidential information of the parties.\nThere are several situations where two parties need to exchange information with one another to determine information that is held in common by both parties. As part of the information exchange, the parties typically also desire that information that is not held in common by the parties (i.e., information that is held by one party but not the other party) be not disclosed to the other party.\nOne example of such a situation occurs in the email marketing business. The use of electronic mail messages (“emails”) as marketing tools has gained widespread popularity in recent years. Emails provide a convenient channel for marketers to reach a large audience of potential consumers. Marketers generate and maintain lists of email addresses and send emails to addresses in the list. The emails may comprise information marketing/advertising vendors' products, hyperlinks (e.g. URLs) to the vendors' websites, and the like. Marketers also commonly run email campaigns targeted against email address lists of another marketer (referred to as the “rental affiliate” or “affiliate”). In such a scenario, the affiliate sends emails on behalf of the marketer to email addresses included in the address list maintained by the affiliate.\nThe recent enactment of laws aimed to reduce spam is requiring marketers to change the manner in which email campaigns are conducted, especially campaigns targeted against address lists of a rental affiliate. For example, under the recently enacted federal anti-SPAM law (“CAN-SPAM”), emails marketers are required to suppress (i.e., not send emails to) unsubscribed email addresses from all email campaigns they initiate. This includes email campaigns implemented by and targeted against the list of a list rental affiliate.\nAccordingly, when a marketer runs an email campaign against a list of email addresses of an affiliate, the new laws require that emails not be sent to email addresses that occur both in the marketer's list of addresses to be suppressed and in the affiliate's list of email addresses. Accordingly, in order to abide by the new laws, it is essential to determine email addresses that occur in the marketer's suppression list and the affiliate's list. These email addresses that are in both lists represent information that is held in common by the two parties, namely the marketer and the affiliate. However, in the process of determining these common addresses, neither the affiliate nor the marketer wants their non-common information (e.g., addresses in the affiliate's list that are not in the marketer's suppression list and email addresses in the marketer's suppression list that are not in the affiliate's list) to be disclosed to the other party.\nHistorically, affiliates have offered suppression services for an additional fee. Typically, a marketer would send a list of email addresses to be suppressed (the “suppression list”) to an affiliate and the affiliate would remove any email addresses found in the suppression list from the affiliate's list prior to conducting the email campaign. However, this approach violates the privacy policies of many marketers, which prohibit them from sharing customer data with third parties such as the affiliate.\nAccording to another technique, the affiliate sends its list of email addresses to the marketer. The marketer then removes email addresses from the affiliate list that are to be suppressed and returns the reduced list of addresses back to the affiliate. The affiliate then conducts email campaigns against only those email addresses in the reduced list received from the marketer. While this approach does not expose the marketer's customer data to the affiliate, it exposes the affiliate's customer data to the marketer. This is generally unacceptable to the affiliate and may violate its privacy policies.\nAccordingly, a challenge in implementing suppression of email addresses between an affiliate and a marketer is that both parties view their data as private. Neither the affiliate nor the marketer wishes to disclose their customer and suppression information to the other party, and yet both have to cooperate to meet the legislative requirements of suppressing unsubscribed addresses."}
{"text": "1. Field of the Invention\nThe present invention relates to a weather radar, and weather observation method configured to three-dimensionally observe a weather phenomenon such as rain, cloud, and the like.\n2. Description of the Related Art\nIn the conventional weather radar of the parabolic antenna type, a thin beam called a pencil beam is formed, the beam is rotated 360° in the horizontal direction to acquire observation data in one plane, and then the antenna elevation angle is increased stepwise to continue to acquire the next one plane data, thereby collecting three-dimensional precipitation data (see, for example, “Revised Radar Technique” supervised by Takashi Yoshida, Institute of Electronics, Information and Communication Engineers, Oct. 1, 1996, first edition, “Chapter 9, Meteorological Radar”, pp. 238-240). In order to carry out this observation sequence, a time of about five to ten minutes is required, and hence temporal/spatial resolution sufficient for observation of a cumulonimbus cloud or the like changing every moment has not been obtained."}
{"text": "Numerous articles of commerce are manufactured by subjecting a moldable composite sheet to a molding process. The moldable composite sheets usually comprise a thermoplastic resin such as a polypropylene resin which is reinforced with reinforcement fibers such as a glass fiber mat. These composite sheets are commonly referred to as glass mat thermoplastics or GMT. These composite sheets can be heated to soften the resin and molded under pressure to make various articles such as parts for automobiles, cabinetry for computers and other electronic equipment, components for small appliances, and recreational, agricultural and other durable consumer and industrial goods. A typical molding process for these composite sheets is compression molding in which the moldable composite sheet is heated, pressed into a mold, and subjected to pressure in a relatively cold mold to form the molded part. The resin and reinforcement material flow during the molding process to fill out the entire volume of the mold.\nA compression molding technique is disclosed in U.S. Pat. No. 4,970,044 to Kim et al., which is hereby incorporated by reference. Kim et al. teach that a deformable plastic film can be placed on the surface of a thermoplastic charge prior to compression molding. The film acts as an insulator to prevent the relatively cold mold from cooling the thermoplastic charge so fast that the surface of the charge becomes moldable. Kim et al. further teach that the film can be compatible with the thermoplastic charge so that the molding process bonds the film to the charge, thereby producing a molded composite product.\nMoldable composite sheets and molded composite products often contain additives that are nonresinous but are used to enhance the molded product. An example of an additive is a colorant. Various pigments are added to the resin so that the molded product has the desired color. One of the problems with using additives in molded composite sheets is the cost of the additives. Frequently, the cost of the additive is a substantial fraction of the total raw material cost of the molded product. It would be advantageous if the cost of molded composite products containing additives could be reduced."}
{"text": "Point of sale (POS) terminals, price scanners, and handheld devices used by sales associates in places like retail establishments often incorporate dedicated bar code readers (or bar code scanners) to optically sense a bar code image and produce an electronic signal corresponding to the sensed image. A variety of dedicated bar code readers exist. Dedicated bar code readers typically include one or more illuminating light sources, an aiming light source, an imaging sensor, and electronics to process and interpret the data received. Generally, these readers can operate quickly and reliably for scanning bar codes on merchandise, but have limited interactive and communication capabilities and can be expensive, particularly when they are of the 2D variety."}
{"text": "A typical example of the conventional art is disclosed in U.S. Pat. No. 4,726,689, in which a plurality of exhaust grooves are provided on a housing for a direct air bearing disclosed therein, with several of the grooves nearest a vacuum chamber being exhausted by the vacuum pump and the groove furthest from the vacuum chamber being opened to the atmosphere.\nHowever, the conventional art described above has several disadvantages.\nThe exhaust system of the conventional art requires a large number of pumps in order to operate properly, which increases the size and weight of the bearing. Moreover, the exhaust system is located outside the vacuum chamber, which is unsuitable for use with a stage capable of two-dimensional movement like the XY stage."}
{"text": "In one type of nuclear reactor, control rods are selectively inserted and withdrawn from a nuclear reactor vessel for controlling the operation thereof. Each of the control rods is typically positioned by a conventional control rod drive which includes a ball screw or spindle threadingly engaging a ball nut for raising and lowering the ball nut as the spindle is rotated either clockwise or counterclockwise. A hollow piston rests upon the ball nut at one end thereof and at its other end is conventionally joined to the control rod. Displacement of the ball nut provides displacement of the piston which in turn inserts or withdraws the control rod in the reactor vessel.\nIn order to achieve faster insertion of the control rod than could be obtained by normal rotation of the ball spindle, which is conventionally referred to as a scram operation, a rapid flow of high-pressure water is injected through the control rod drive past the piston for lifting the piston off the ball nut in a relatively short time for quickly inserting the control rod into the reactor vessel. The high-pressure water is channeled to the control rod drive through a scram line pipe attached to a high-pressure water accumulator.\nIn one type of occurrence which allows for rapid backflow of the water past the piston, due to, for example, a break in the scram line, the backflow may cause a large reverse pressure on the piston which in turn provides a back force on the control rod ball nut. This back force can cause reverse rotation of the ball spindle with corresponding withdrawal of the control rod. Withdrawal of one of the control rods due to such a backflow occurrence may cause damage to adjacent fuel in the reactor vessel, requiring replacement thereof leading to undesirable down time of the reactor and economic losses.\nIn order to prevent the above occurrence, a conventional electromechanical brake is provided in the control rod drive for holding the ball spindle from rotating unless the brake is energized. The brake is sized for restraining rotation of the ball spindle against such forces due to backflow of water over the piston when the control rod drive motor is not operating. And, when the control rod drive motor is operating, the motor itself is sized for providing adequate torque for resisting the forces due to the backflow of water in the event of the above-described occurrence.\nTo ensure operability of the brake, the brake is periodically tested. However, the brakes are located adjacent to the reactor vessel, which is inaccessible during operation of the reactor due to the radiation field emanating from the reactor vessel. The radiation field continues at reduced levels also during shutdown of the reactor, which would require inspectors to wear suitable protective clothing and limit their time in the area. In one nuclear reactor embodiment, there are about 205 control rod drives, including a respective number of brakes, which would necessarily require a substantial amount of time for testing all of the brakes. Testing of the brakes during reactor shutdown would, therefore, be relatively costly to accomplish, which is additionally economically undesirable since the reactor is not operating for producing power.\nSince conventional electromechanical brakes typically utilize braking pads for restraining rotation of a rotor disc, they are subject to slippage. Slippage can result in undesirable partial withdrawal of the control rod during backflow ocurrence, and also requires additional means for effectively testing the torque-resisting capability of the brake."}
{"text": "Azithromycin, 9-deoxo-9a-aza-9a-methyl-9a-homoerythromycin A (N-methyl-11-aza-10-deoxo-10-dehydroerythromycin A: IUPAC) of formula (II) disclosed in U.S. Pat. Nos. 4,517,358 and 4,474,768, is an azalide-type semi-synthetic macrolide antibiotic, useful for treating bronchial infection, sexual contact infection and dermatological infection (See Kirste and Sides; Antimicrob. Agents Chemother., 33, 1419(1989)).\n\nAzithromycin is known to exist in three forms, the anhydride, monohydrate and dihydrate forms. These forms have been identified by powder X-ray diffraction and differential scanning calorimetric studies.\nAzithromycin anhydride, which is disclosed in U.S. Pat. No. 4,517,359, is non-crystalline product and thus, its highly hygroscopic property is not suitable for pharmaceutical formulation.\nFurther, azithromycin monohydrate (mp. 136° C.), as described in U.S. Pat. No. 4,474,768 and WO Publication No. 89/00576, is crystalline but it has also hygroscopic property, making it difficult to maintain its water content at a constant level.\nWO Publication No. 89/00576 discloses a process for preparing azithromycin dihydrate (mp. 126° C.) from azithromycin monohydrate by recrystallizing from a mixture of tetrahydrofuran, water and a C5˜C7 aliphatic hydrocarbon. Although the dihydrate is less hydroscopic than the monohydrate, the water content thereof must be carefully maintained during a vaccum drying step at a relatively low temperature. Such a water content controlling procedure is, however, not sufficient for removing the toxic aliphatic hydrocarbon solvent rigorously used in the recrystallization procedure. On the other hand, vacuum drying in higher temperature may result in formation of azithromycin dihydrate having undesirable water content.\nAccordingly, many attempts have been made to develop a novel crystal or solvate form of azithromycin. For example, EP Publication No. 0,984,020 discloses a clathrate of azithromycin monohydrate with isopropanol of formula (III).\n\nWO Publication No. 00/32203 discloses an ethanol solvate of azithromycin hydrate of formula (IV).\n\nHowever, there has existed a need to develop an improved crystal form of azithromycin crystal suitable for pharmaceutical applications."}
{"text": "Such devices are known in the prior art. In particular, portable electronic devices are known displaying an altimeter function, usually intended to be used when hiking.\nHowever, in order to practise some activities, an indication of the altitude is not sufficient for the wearer of the device. In particular, in the practice of aeronautical sports involving piloting unpowered aircraft, it is important for the pilot to have information relative to his vertical speed of movement. The onboard instruments existing on transport aircraft present a large number of functions which are not all provided for the pilots of light aircraft of glider, microlite, para-glider or delta plane type.\nFurthermore, in the practice of some other sports, such as parachute jumping, hiking or even mountain biking, the fact that the sportsman has available information about his altitude is useful and even essential in that which relates to free fall. However, it can be of interest to persons performing this type of activity to have complementary information available."}
{"text": "Computer security is a technology associated with information security of computer networks. The objective of computer security includes protection of information and property from theft, corruption, or natural disaster, while allowing the information and property to remain accessible and productive to its intended users."}
{"text": "In conventional systems, a viewing experience of a user is limited to the images displayed within the boundaries of a television or computer screen. Thus, an area outside of the television or computer screen can be found distracting to a user, as it is not pertinent to the images on the screen."}
{"text": "As semiconductor designs continue marching towards smaller geometry, back end of line (BEOL) copper wire formation with lower dielectric constant material becomes a significant challenge. Damascene integration techniques are commonly used to form copper wires and vias. Conventional damascene integration methods involve forming trenches in an insulator layer, depositing a thin diffusion barrier layer and seed layer on surfaces of the trenches, electroplating to fill the trenches in the insulator layer, and planarizing, e.g., with chemical mechanical polish (CMP). As geometry becomes increasingly smaller, e.g., sub-micron, it is very difficult to sputter and plate copper material into tiny damascene trenches. Moreover, a significant amount of copper (Cu) wire thickness variation exists when Cu-CMP is employed."}
{"text": "Reducing the cost of aerospace parts requires a switch to reliable, repeatable processes that yield parts, which meet the rigorous requirements of aerospace. In aerospace, the issues of absolute weight, shape, strength, and durability must be traded during the design of a part. A factor often ignored during this design trade is the cost of manufacture. Often, however, safety forces part designs that are expensive to manufacture. Resin transfer molding (RTM) is promising as a process to replace complicated built-up metal assemblies with one-piece fiber-reinforced resin composite alternatives. One challenge with RTM is ensuring that the resin-to-fiber ratio in the molded part is adequate everywhere in the part so that the molded part will have the design strength and durability. The present invention is an improvement in RTM using a seal pin in the mold die to index the preform for the part reliably and precisely on the mold cavity on an inner mandrel that is completely sealed within the die set.\nResin transfer molding (RTM) is a process by which a resin is pumped at low viscosities and low pressures into a closed mold die set containing a preform of dry fabric to infuse resin into the preform and to make a fiber-reinforced composite part. The RTM process can be used to produce at low cost composite parts that are complex in shape. These parts typically require continuous fiber reinforcement along with inside mold line and outside mold line controlled surfaces. The ability to include and place continuous fiber reinforcement in large and small structures sets RTM apart from other liquid molding processes.\nSome RTM mold die sets require outer dies to define a cavity and an inner mandrel that seats in the cavity to define a hollow region in the molded part. Often the mandrel must be assembled from multiple parts so that it can be disassembled and removed from the molded part. RTM die sets that used mandrels usually are configured so that the mandrel is inserted through an opening in the dies. Sealing the base plate of the mandrel with the opening in the dies was challenging. With the seal/index pins of the present invention, the mandrel can be positioned entirely within the mold cavity, allowing a simpler seal system between the dies."}
{"text": "The invention related to new method for electronic card games, such as poker."}
{"text": "EPROMs and Flash E.sup.2 PROMs (hereafter collectively, PROMs) have several structures which allow them to hold a charge without refresh for extended periods of time (see FIG. 1). The charge itself is stored on a \"floating gate\" 10 also referred to as Poly 1 or P1, which is a structure of polycrystalline silicon (hereafter, poly) surrounded on all sides by a layer of oxide 12. Located superjacent and parallel to this P1 structure is another poly structure, the word line or \"control gate\" 14 or P2. P1 10 and P2 14 act as the two plates of a capacitor. Below the P1 layer are two N+ junctions, one which acts as the transistor source 16 and the other as the drain 18, which are doped into a p-type substrate 20. The portion of the substrate 20 between the source 16 and the drain 18 is the channel 22. The cell of FIG. 1 functions like an enhancement-type n-channel metal oxide semiconductor field effect transistor (MOSFET) with two gates of poly.\nThere are many ways to program a flash E.sup.2 PROM. For example, a high voltage such as 12 V, is applied to the control gate. Then a voltage pulse such as 8 V is applied between source and drain. The large positive voltage on the control gate establishes an electric field in the insulating oxide. This electric field attracts the electrons generated from the so-called avalanche breakdown of the transistor due to the high drain and control gate voltages, and accelerates them toward the floating gate, which they enter through the oxide. In this way the floating gate is charged, and the charge that accumulates on it becomes trapped.\nTo return the floating gate from a charged state to a state with no charge, the charge is returned to the substrate or other electrodes. In an EPROM, this is accomplished with ultraviolet light which excites the electrons past a certain energy state, thereby allowing them to pass through the oxide. In an E.sup.2 PROM, this is accomplished with an electrical field.\nThe voltage which must be applied on the control gate to turn on the transistor is much higher in a device storing a charge than in a device which does not have a voltage potential stored on P1. To read the content of the floating gate, a voltage somewhere between the low and high voltage values (i.e. the threshold voltage V.sub.t) is applied to the control gate. A cell that trips at V.sub.t has no charge stored on P1, while a cell which does not trip is determined to be storing a charge.\nThere are structures which make up a PROM array which are common to several transistors in the array. FIG. 2 shows a top view of an array showing transistor sources 16, drains 18, digit lines 24, floating gates 10, and control lines 26 which form control gates 14 as they pass over the floating gates 10. Also shown as a dotted line is the \"active area\" 28 interspersed with areas of field oxide 30. As shown in FIG. 2, a single control line 26 is common to all transistors in a single column, and when selected it activates all transistors in the column. The source regions 16, which run parallel with the control lines 26, are common to all transistors in two adjacent columns. Individual transistor drains 18 are common to two transistors in adjacent columns. The digit (or bit) lines 24 are common with the drains 18 of all transistors in a single row.\nTo read the datum on a floating gate 10, the control line 26 of the cell to be read is activated which causes all transistors in the selected column to become active and to output the cell information on their respective digit lines 24. The information on the digit line 24 which corresponds to the cell to be read is obtained with a sense amplifier (not shown), with one sense amp for each digit line.\nThe active area 28 is defined during the manufacturing process after the thin pad oxide is initially formed on the surface of the wafer. After the pad oxide is formed on the surface of the wafer a layer nitride is formed and a patterned layer of photoresist is formed over the nitride. The exposed nitride is removed, and the exposed oxide is further oxidized. The nitride prevents the growth of the oxide underneath it, while allowing the exposed oxide to grow. The area under the nitride which remains thin is referred to as the gate oxide and is also called the active area. The exposed oxide which thickens becomes the field oxide. During a spacer etch, the gate oxide will erode away to bare silicon, as it is thinner than the field oxide. The bare areas can then be doped, with the field oxide protecting the other silicon areas from being doped.\nThere are many methods used to manufacture Flash E.sup.2 PROM and EPROM memory, one of which is shown in FIGS. 3-7. To obtain the structure of FIG. 1, a first blanket layer of oxide, the pad oxide, is formed on top of the silicon substrate. Nitride is patterned in a criss-cross pattern on the oxide to define the active area 28 as shown in FIG. 3, and the field oxide 30 is formed from the exposed pad oxide while the unexposed pad oxide becomes the gate oxide 28. As shown in FIGS. 4A and 4B, a first blanket layer of poly 40 which will make up P1 is formed on the appropriate areas of the gate oxide 28, and a second layer of oxide 42 and a layer of nitride 44 (or a nitride-oxide sandwich), which separate P1 from a second poly layer is formed on the first poly layer 40. The P1 layer is patterned with a layer of resist 46 in rows, perpendicular to the source lines of the active area. As shown in FIG. 5, after a P1 sidewall oxidation (not shown) a second blanket poly layer 50, P2, is formed on the nitride layer 44 (or nitride-oxide), and an optional oxide layer 52 is formed on the Poly 2 50 layer, and photoresist 54 is patterned on the surface in columns perpendicular to the rows of P1 40. An etch forms the floating gates 10 and control gates 14. Referring to FIG. 6, a third blanket layer of oxide 60 is formed and etched, which forms the spacers 70 as shown in FIG. 7. During the spacer etch, exposed gate oxide 28 is also removed, thereby exposing the areas of silicon which will later become the transistor diffusion areas.\nAfter the structure of FIG. 7 is formed, the substrate 20 is doped to form the N+ areas of the transistor sources 16 and drains 18, and conductive areas which couple the sources between the rows. An implant source-drain can also be formed before the spacer formation for additional overlap of the impurity regions under the spacers.\nA second method of forming the Flash E.sup.2 PROM structure is described in the article \"A 5-Volt Contactless Array 256KBIT Flash EEPROM Technology,\" M. Gill et al, IEDM, 1988, pp 428. This structure, which is formed with a buried N+ line process, has higher digit line capacitance, more process complexity, and also has a high degree of lateral diffusion in the buried digit line and is not easily scalable.\nWith the advent of laptop computers, reducing the power consumption of devices has become a major design focus for engineers. The high voltage required to program and operate Flash EPROMS can decrease the usable battery life of portable and notebook computers.\nA Flash EPROM design is described in \"Buried Source-Side Injection (BSSI) for Flash EPROM Programming,\" Cetin Kaya et al, IEEE, 1992, pp. 465-467, which is incorporated herein by reference. Kaya, et al. claims the cell, which has a small cell area, can be programmed in a voltage range of 3-4 V. The device structure comprises the use of a buried source-side injector which requires the formation of a source junction 0.1-0.2 microns (.mu.) under the silicon substrate surface, and a conventionally designed drain junction. To program the cell the source is grounded while 13 V and 3.5 V are applied to the gate and to the drain respectively. To manufacture the buried source injector, a self-aligned high-energy n-type dopant implant is performed at the source side, which is followed by a heavy dose of As at the drain side. Subsequently, after the sidewall oxide is formed, additional impurities are deposited in the source to reduce the resistance of the diffusion and to form a low-resistance contact region.\nIn effect, the structure that allows the low-voltage program in the cell of Kaya, et al. is the \"sharp corner\" of impurities at the source close to the floating gate which forms the source injector. Maximizing the sharpness of the implanted region is a goal of the high-energy implant, but the implant is difficult to control and the shape of the injector will vary widely depending on the subsequent process heat treatment. A cell which has the advantages of the Kaya et al. cell but which allows the sharpness of the source injector to be maximized would be a desirable structure."}
{"text": "In most mobile communication systems of today, there are specific requirements regarding synchronization of a base station and a mobile terminal in order to secure a correct data transmission. Examples of such systems are the Universal Terrestrial Radio Access (UTRA) and Evolved UTRA.\nIn Evolved UTRA, Single-Carrier Frequency Division Multiple Access (SC-FDMA) may be used as multiple access scheme for the uplink communication. The transmission scheme of SC-FDMA is the so-called Discrete Fourier Transform-spread Orthogonal Frequency Domain Multiplexing (DFT-spread OFDM), which can be seen as OFDM with pre-coding. Whereas OFDM, which produces a multi-carrier signal, has a high peak-to-average ratio (PAPR), the DFT pre-coding gives a single-carrier signal with lower PAPR. The low PAPR serves to extend the coverage and to reduce the battery drain in the mobile.\nIn DFT-spread OFDM, cyclic prefix is used to achieve equalization in the frequency domain. However, a requirement for successful equalization in DFT-spread OFDM, as well as in OFDM, is that the signals transmitted from all mobile terminals in the cell are synchronized in such a manner that the delay spread of the signal plus the spread in the time of arrival is less than the duration of the cyclic prefix. Therefore, it is required that each transmitting mobile terminal is synchronized to within a fraction of the duration of the cyclic prefix before it can transmit data.\nIn Evolved UTRA, synchronization is performed both in uplink and downlink. In one step of the synchronization, downlink synchronization, the mobile terminal synchronizes (or locks) to the carrier frequency and the frame timing of the base station. This synchronization, however, is not sufficient to ensure that the base station can properly receive the signals from the mobile terminal, since mobile terminals may be located at various distances relative to the base station.\nConsequently, further synchronization, uplink synchronization, is needed since the distance between a base station and a mobile terminal, and hence the round trip time, is in general unknown.\nIn Evolved UTRA, a random access channel (RACH) supports uplink synchronization of the mobile terminals. RACH in Evolved UTRA is contention-based, i.e. any mobile terminal within the cell may transmit on the resource allocated to RACH. Consequently, several mobile terminals may attempt to transmit synchronization signals simultaneously, and in order to reduce the risk that the base station fails to distinguish signals from different mobile terminals, a set of signature sequences is provided, wherein each mobile terminal randomly selects one signature sequence.\nIn UTRA and Evolved UTRA a binary pseudo-random sequence generated by a shift register is modulated by 16-bit Hadamard sequences to produce these signature sequences. Even though these signature sequences in many instances provide good correlation properties, there still exists a need for enhanced detection capabilities to detect a specific signature in presence of other simultaneous signatures, especially at low SIR values."}
{"text": "The invention relates to an alert button. More particularly, the invention relates to a miniature button which may be affixed to clothing and clothing accessories, and is capable of being activated by the wearer to transmit an alert signal to another location.\nGreat advances have taken place in the field of personal communications in recent years. Cellular telephones are now small enough to fit in the shirt pocket. Pagers are capable of receiving not only numeric messages, but also alphanumeric and voice messages. Some wristwatches now incorporate paging devices.\nThese devices provide great versatility in terms of the type of information that may be conveyed. They also demonstrate the extent of current miniaturization technology, wherein a complicated computerized transceiver may now be contained in a very small housing.\nHowever, these devices fall short in terms of providing a simple alert signal in the event of an emergency. A typical cellular phone is unsuitable for providing an emergency alert.\nFirst, even though cell phones have become quite small, they are still not so small that they can be carried everywhere. A cell phone is typically carried in a pocket, or in a purse. However, when one is not carrying a purse, or is wearing clothing without pockets, they will typically not carry a cell phone.\nSecond, it takes time to actually make an emergency communication using a cell phone. The phone must be switched on, several keys must be pressed to dial an emergency number, it can take several seconds for the call to connect, and then one must speak to communicate the emergency. This process can be much too time consuming in certain emergency situations.\nFurther, a cellular telephone is not easy to operate, especially for small children, the elderly, and people with physical or mental disabilities. All of the steps necessary to operate the phone may be well beyond one of these people. Further, it is not always desirable to give a cellular telephone to a child or a teenager, who might use the telephone indiscriminately in non-emergency situations.\nPagers are not at all suitable for communicating the presence of an emergency situation. Pagers are only capable of receiving messages--not transmitting them.\nU.S. Pat. No. 4,899,135 discloses a child monitoring device, in which an alert signal is generated when a child, wearing a transmitter, strays from the parent beyond a predetermined threshold distance. U.S. Pat. No. 4,785,291 to Hawthorne also discloses a distance monitor for child surveillance.\nU.S. Pat. No. 5,337,041 to Friedman discloses a personal safety guard system for a stray person or pet. The device is typically worn by a pet, and can alert the pet when their owner is looking for them. In addition, the device can emit an audible alarm signal from the pet to notify passersby that the pet is lost.\nU.S. Pat. No. 5,021,794 to Lawrence, discloses a personal emergency locator system which remotely activates and \"homes in\" on a miniaturized transceiver which is worn upon a person, such as a lost child.\nWhile these units may be suitable for the particular purpose employed, or for general use, they would not be as suitable for the purposes of the present invention as disclosed hereafter."}
{"text": "Conventionally, monitors enabling an observer to view two-parallax images captured from two viewpoints stereoscopically by using a specific device, such as a pair of stereoscopic vision glasses, have been in practical use. Furthermore, in recent years, monitors enabling an observer to view multi-parallax images (e.g., nine-parallax images) captured from a plurality of viewpoints stereoscopically with the naked eyes by using a beam control element, such as a lenticular lens, have also been in practical use. Such two-parallax images and nine-parallax images displayed on monitors enabling stereoscopic vision may be generated by estimating depth information of an image captured from one viewpoint and performing image processing with the information thus estimated.\nAs for medical image diagnosis apparatuses, such as X-ray computed tomography (CT) apparatuses, magnetic resonance imaging (MRI) apparatuses, and ultrasound diagnosis apparatuses, apparatuses capable of generating three-dimensional medical image data (hereinafter, referred to as volume data) have been in practical use. Conventionally, volume data generated by such a medical image diagnosis apparatus is converted into a two-dimensional image by various types of image processing, and is displayed two-dimensionally on a general-purpose monitor. For example, volume data generated by a medical image diagnosis apparatus is converted into a two-dimensional image that reflects three-dimensional information by volume rendering processing, and is displayed two-dimensionally on a general-purpose monitor.\nIn the conventional technology, however, when sequential multi-parallax images are generated and displayed in real time, processing load may increase, thereby making it difficult to display the sequential multi-parallax images smoothly."}
{"text": "1. Field of the Invention\nThe art disclosed by the present specification relates to a vehicle mounting structure for an electronic device in an engine room or a motor room.\n2. Description of Related Art\nAn engine room or motor room of a vehicle is mounted with various devices including an engine or a running motor. In recent years in particular, electric vehicles including hybrid vehicles have been in widespread use, and these vehicles have been mounted with various electronic devices. An inverter that generates an alternating electric current for driving the running motor is a typical electronic device mounted on a vehicle. Various ingenuities have been proposed as to the mounting of important electronic devices, especially inverters.\nFor example, in Japanese Patent Application Publication No. 2004-328870 (JP-2004-328870 A), there is disclosed a structure in which an inverter is fixed above a motor with a gap therebetween. Besides, in Japanese Patent Application Publication No. 2008-248936 (JP-2008-248936 A) and Japanese Patent Application Publication No. 2011-115013 (JP-2011-115013 A), there is disclosed a structure in which an inverter is fixed to a vehicle body via a rubber vibration isolator. A mounting metal fitting for fixing an electronic device to a vehicle body or another component is often referred to as a bracket. In the present specification as well, a metal fitting for mounting an electronic device is referred to as a bracket. As disclosed in Japanese Patent Application Publication No. 2008-248936 (JP-2008-248936 A) and Japanese Patent Application Publication No. 2011-115013 (JP-2011-115013 A), a bracket and a rubber vibration isolator are often used as a set. Besides, in Japanese Patent Application Publication No. 8-290720 (JP-8-290720 A), a set of a bracket and a rubber vibration isolator is used to fix a motor to a vehicle body. The bracket of Japanese Patent. Application Publication No. 8-290720 (JP-8-290720 A) is provided with a flange for increasing the strength thereof."}
{"text": "Golf clubs take various forms, for example a wood, a hybrid, an iron, a wedge, or a putter, and these clubs generally differ in head shape and design (e.g., the difference between a wood and an iron), club head material(s), shaft material(s), club length, and club loft.\nWoods and hybrids generally have a longer shaft and lower loft than irons and wedges. Thus, a golf ball that is struck with a wood or a hybrid generally travels a greater distance than a golf ball struck with an iron or a wedge. While a longer shaft and a lower loft provide increased golf ball travel distance, this combination also results in less forgiveness. The longer shaft requires a golfer to stand farther away from the golf ball at address. This leads to greater difficulty during the golf swing to return the club head squarely to impact the golf ball. A golf club that is slightly open or slightly closed at impact results in reduced accuracy as the golf ball is not launched on the desired target line. Further, the higher swing speeds from the longer length shaft can lead to greater difficulty in making consistent contact with the center or “sweet spot” of the golf club face. Off-center contact can lead to imparting increased side spin on the golf ball. At reduced lofts of woods and hybrids, less back spin is imparted on the golf ball at impact, further exacerbating imparted side spin and leading to undesirable hooks or slices, which further decrease accuracy.\nTo improve directional forgiveness, golf club manufacturers have made efforts to increase the moment of inertia of a golf club at impact. The moment of inertia (or “MOI”) is a measure of a body's resistance to angular acceleration, or twisting. The higher the MOI of a golf club head, the more the golf club head resists twisting at impact, improving golf ball accuracy, especially on off-center contact (or mishits). In addition, the increased stability of a higher MOI golf club head results in a golf ball losing less ball speed on off-center contact due to reduced energy loss associated with reduced twisting. A higher MOI of a golf club head further increases consistency in spin rate and launch angle of a golf ball on off-center contact.\nWhile woods and hybrids have a variety of known designs, there is a need for enhancing directional forgiveness (e.g., a reduction in side-to-side variation) to improve accuracy, especially on off-center hits (e.g., contact of the golf ball with a location on the golf club face other than the sweet spot)."}
{"text": "Paper, as that term is used herein, refers to the wide variety of cellulose-based sheet materials used for writing, printing, packaging, and other applications. Paper may be used, for example, but without limitation, in the following applications: as paper money, bank notes, stock and bond certificates, checks, and the like; in books, magazines, newspapers, and art; for packaging, e.g., paper board, corrugated cardboard, paper bags, envelopes, wrapping tissue, boxes; in household products such as toilet paper, tissues, paper towels and paper napkins; paper honeycomb, used as a core material in composite materials; building materials; construction paper; disposable clothing; and in various industrial uses including emery paper, sandpaper, blotting paper, litmus paper, universal indicator paper, paper chromatography, battery separators, and capacitor dielectrics.\nPaper is generally produced by pulping a cellulosic material to form a pulp containing cellulosic fibers, amalgamating the cellulosic fibers to form a wet web, and drying the wet web. In the finished paper, the fibers are held together by mechanical interlocking and hydrogen bonding. Pulping may be accomplished in a number of ways, for example: using a chemical process (e.g., the Kraft process), a mechanical process (groundwood), or thermomechanical process (TMP). The amalgamating and drying steps are generally performed using a high speed paper machine.\nThe most common source of cellulosic fibers is wood pulp from trees. Pulp is also derived from recovered (“recycled”) paper. Vegetable fiber materials, such as cotton, hemp, linen, and rice, are also used. Other non-wood fiber sources include, but are not limited to, sugarcane, bagasse, straw, bamboo, kenaf, jute, flax, and cotton. A wide variety of synthetic fibers, such as polypropylene and polyethylene, as well as other ingredients such as inorganic fillers, may be incorporated into paper as a means for imparting desirable physical properties.\nFor many applications, it is desirable that paper have high strength and tear resistance, even in very thin sheets, for example, when the paper is used in packaging, in industrial applications, as money, and in other applications that require strength and durability. It is also generally desirable that paper exhibit good printability characteristics, with the particular characteristics depending to some extent on the printing process in which the paper will be used."}
{"text": "1. Field of the Invention\nThis invention relates to a process for operating a charge coupled arrangement in a two-phase operation, in which an insulating layer is provided on a substrate of semiconductor material and contains electrodes at a first level spaced apart therein and insulated from the semiconductor substrate, and second electrodes at a second level arranged above the interspaces between the electrodes of the first level and electrically insulated from the electrodes of the first level, and wherein each electrode of the second level is electrically connected to and adjacent electrode of the first level and a zone, doped opposite to that of the semiconductor substrate, is provided in the substrate beneath the electrodes and beneath the insulating layer.\n2. Description of the Prior Art\nCharge-coupled arrangements are well known in the art. For example, the publication \"The new concept for memory and imaging: charge coupling\" by Laurence Altman, appearing in the June 21, 1971 issue of the periodical Electronics at Pages 50 to 58, describes a SCCD charge-coupled arrangement (surface charge-coupled device) for two-phase operation in an aluminum-silicon-gate technique. In order to change the transmission direction, it is necessary to change the pulse program for the pulse train lines. However, this means that in each case each second aluminum electrode and each second silicon electrode must be connected to an individual drive electrode and that these four lines must be accessible from the exterior of the device. In the event that adjacent aluminum and silicon electrodes are permanently connected, it is not possible, however, to change direction in this manner.\nThe \"Bell System Technical Journal\", Vol. 54, September 1972, No. 7, Pages 1635-1640 describes a BCCD (buried channel charge-coupled device) arrangement. In such an arrangement, below the electrodes between the semiconductor substrate and the insulating layer is arranged a zone which is oppositely doped in respect of the doping of the substrate."}
{"text": "In general, a conventional vehicle navigation apparatus is designed to navigate from a current position to a destination set up by a user. Meanwhile, the conventional vehicle navigation apparatus has been known to be constructed to record, on a recording medium, information representing a number of times and conditions of arrival to respective past arrived positions for the purpose of setting priorities for the respective past arrived positions on the basis of the recorded information, and estimating a destination from the past arrived positions on the basis of the set priorities, even if a destination is not set (see for example Patent Document 1)."}
{"text": "The present invention relates to converting a fuel to combustible gas.\nEuropean patent application No. 80300003.3, publication No. 0013590, describes and claims a method of converting a fuel (which may be a liquid fuel or a solid fuel or a mixture of liquid and solid fuels) to a substantially inerts-free combustible reducing and/or synthesis gas product, the method comprising passing the fuel into a fuel conversion zone comprising a conversion bed of solid particles which are maintained at an elevated fuel conversion temperature and fluidized by a fluidizing gas which is passed into the bottom of the conversion bed and rises upwardly therethrough, the conversion bed comprising particles containing reactive calcium sulfate (CaSO.sub.4) which constitutes substantially the sole source of oxygen used in the conversion of the fuel, and the fluidizing gas passed into the bottom of the conversion bed comprising at least one gas phase and/or vapour phase substance which serves to promote and/or mediate the transfer of oxygen to the fuel and/or partially converted derivatives thereof from the CaSO.sub.4 (which is thereby reduced to calcium sulfide, CaS), the fluidizing gas being substantially free of (a) molecular oxygen, and (b) unreactive and/or inert materials which would dilute the combustible gas product.\nIn preferred embodiments, particles from the conversion zone are fluidized by molecular oxygen-containing gas in an oxidation zone separate from the fuel conversion zone so as to oxidize calcium sulfide in said particles to CaSO.sub.4, and particles which have been fluidized in the oxidation zone are returned to the fuel conversion bed for use in converting further quantities of fuel, the molecular oxygen-containing gas being in excess to substantially suppress the empirical reaction: EQU 2CaS+30.sub.2 .fwdarw.2CaO+2SO.sub.2 ( 1)\nIn preferred embodiments, the conversion bed contains particles comprising calcium oxide (CaO) to react with and fix sulfur from the fuel as CaS so that the combustible gas product has a low sulfur content, and particles from the conversion zone are preferably fluidized by a gas containing molecular oxygen in a regeneration zone so as to convert CaS to CaO with the liberation of sulfur as sulfur oxide (principally SO.sub.2) and/or elemental sulfur, and particles containing CaO are returned to the conversion bed directly and/or indirectly for use in reacting with and fixing sulfur from further amounts of fuel undergoing conversion.\nMany multi-step processes, such as the fuel-conversion process as hereinabove referred to, employ a discrete vessel for each step. A drawback of this almost universal practice is that the flow of particles from one vessel to another is through pipes or conduits which for practical reasons must be so positioned that the flow pattern of particles within each vessel departs appreciably from the ideal."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of wind turbines that use environmental air current and winds to generate electrical energy with an electrical generator.\n2. Background of the Art\nThe wind turbine is a device that converts kinetic energy from the wind into mechanical energy. If the mechanical energy is used to produce electricity, the device may be called a wind generator or wind charger. If the mechanical energy is used to drive machinery, such as for grinding grain or pumping water, the device is called a windmill or wind pump. Developed for over a millennium, today's wind turbines are manufactured in a range of vertical and horizontal axis types. The smallest turbines are used for applications such as battery charging or auxiliary power on sailing boats; while large grid-connected arrays of turbines are becoming an increasingly large source of commercial electric power. Wind turbines can rotate about either a horizontal or a vertical axis, the former being both older and more common horizontal axes.\nHorizontal axis—These are defined as components of a horizontal axis wind turbine (gearbox, rotor shaft and brake assembly) after being lifted into position. Horizontal is nominally in reference to perpendicularity to gravity and/or parallelism with the horizon. Those positions are approximate, and with devices that may also shift with respect to the horizon to adjust for wind speed, that approximation may be significant.\nHorizontal-axis wind turbines (HAWT) have the main rotor shaft and electrical generator at the top of a tower, and must be pointed into the wind. Small turbines are pointed by a simple wind vane, while large turbines generally use a wind sensor coupled with a servo motor. Most have a gearbox, which turns the slow rotation of the blades into a quicker rotation that is more suitable to drive an electrical generator.\nSince a tower produces turbulence it, the turbine is usually positioned upwind of its supporting tower. Turbine blades are made stiff to prevent the blades from being pushed into the tower by high winds. Additionally, the blades are placed a considerable distance in front of the tower and are sometimes tilted forward into the wind a small amount.\nDownwind machines have been built, despite the problem of turbulence (mast wake), because they don't need an additional mechanism for keeping them in line with the wind, and because in high winds the blades can be allowed to bend which reduces their swept area and thus their wind resistance. Since cyclical (that is repetitive) turbulence may lead to fatigue failures, most HAWTs are of upwind design.\nTurbines used in wind farms for commercial production of electric power are usually three-bladed and pointed into the wind by computer-controlled motors. These have high tip speeds of over 320 kilometers per hour (200 mph), high efficiency, and low torque ripple, which contribute to good reliability. The blades are usually colored light gray to blend in with the clouds and range in length from 20 to 40 meters (66 to 130 ft) or more. The tubular steel towers range from 60 to 90 meters (200 to 300 ft) tall. The blades rotate at 10-22 revolutions per minute. At 22 rotations per minute the tip speed exceeds 300 feet per second (91 m/s). A gear box is commonly used for stepping up the speed of the generator, although designs may also use direct drive of an annular generator. Some models operate at constant speed, but more energy can be collected by variable-speed turbines which use a solid-state power converter to interface to the transmission system. All turbines are equipped with protective features to avoid damage at high wind speeds, by feathering the blades into the wind which ceases their rotation, supplemented by brakes.\nWind turbines are designed to exploit the wind energy that exists at a location. Aerodynamic modeling is used to determine the optimum tower height, control systems, number of blades and blade shape. Wind turbines convert wind energy to electricity for distribution and typically are not combined with any storage capability. Conventional horizontal axis turbines can be divided into three components.                The rotor component, which is approximately 20% of the wind turbine cost, includes the blades for converting wind energy to low speed rotational energy.        The generator component, which is approximately 34% of the wind turbine cost, includes the electrical generator, the control electronics, and most likely a gearbox (e.g., planetary gearbox, adjustable-speed drive or continuously variable transmission) component for converting the low speed incoming rotation to high speed rotation suitable for generating electricity.        The structural support component, which is approximately 15% of the wind turbine cost, includes the tower and rotor yaw mechanism.        \nA 1.5 MW wind turbine of a type frequently seen in the United States has a tower 80 meters high. The rotor assembly (blades and hub) weighs 48,000 pounds (22,000 kg). The nacelle, which contains the generator component, weighs 115,000 pounds (52,000 kg). The concrete base for the tower is constructed using 58,000 pounds (26,000 kg) of reinforcing steel and contains 250 cubic yards (190 cubic meters) of concrete. The base is 50 feet (15 m) in diameter and 8 feet (2.4 m) thick near the center.\nAbove 10 m/s, where power regulation caused blade coning, the early data recorded was erratic and use of these figures would distort further data processing. A cut off point was set at 3500 seconds and the readings taken before this point were not used for the detailed calculations. As the wind speed during the day did not fall below 6 m/s, and regulation cut in at 10 m/s, the final results were based on a narrow band of rotor speeds between 3.1 and 4.4 rev/s\nThe universal standard to assess the performance of a wind turbine is a plot of the Power Coefficient versus the Tip Speed Ratio performance curve. This curve tells you how efficiently a turbine converts the energy in the wind to electricity. Various methods were used to plot this curve. The results are show in Figure."}
{"text": "Stringed instruments of the lute type, such as all types of guitars, bass guitars, and other stringed instruments in which the strings are supported by a neck, are characterized by a first hand regulating the sounding length of the strings, and the other hand plucking or bowing the strings to produce a sound, with the pitch adjusted by the established sounding length of the strings.\nThe movement and force needed to establish chords or tones with the first hand are swift, and may need significant application of force. During play, musicians may experience various degrees of discomfort and fatigue of the hands and fingers, and may suffer temporary or chronic localized numbness, for example of the entire area of the thumb that comes into contact with the neck of the stringed instrument. In some cases, debilitating injuries may manifest, such as related to tendonitis, carpal tunnel syndrome, nerve compression injuries and various forms of repetitive strain injuries. Such injuries can be sufficiently serious that they can halt or entirely terminate the career of a professional musician.\nVarious warm-up, preparation, and playing techniques have been proposed in the past, and are in use currently, to prevent such injuries, and a number of well-known treatment methods can be used once an injury has manifested. However, up until now, there are no known devices and associated methods available that can directly prevent, reduce, or alleviate injuries while playing a stringed instrument.\nAs such, considering the foregoing, it may be appreciated that there continues to be a need for novel and improved devices and methods for treatment of trigger points."}
{"text": "This invention relates to a drop-on-demand ink jet printing device commonly used for large character printing. In such devices, used for marking products, the on/off control of ink to a nozzle is accomplished by using a separate valve, remote from the orifice, which is connected thereto by a tubing assembly. A typical large character printhead will have a plurality of such valves and nozzles forming, for example, a seven by one printhead which can mark alphanumeric characters on products which pass the printhead.\nThe subsystem defined by the on/off valve, the nozzle and the tubing which interconnects the two has a natural frequency. In typical printheads of this type, the resonant frequency of the system is low. As the valve operating frequency approaches the resonant frequency of the ink delivery system, print quality deteriorates. Thus, even if the valves are capable of operating at higher frequencies, the resonant interference from the ink delivery subsystem reduces maximum printing frequency.\nFor example, current drop on demand printers of the type described operate at about 240 Hz. Currently available valve designs permit theoretical operation to at least 700 Hz. As print speed increases, however, interfering resonant frequencies cause severe print quality deterioration preventing operation at such higher frequencies.\nIt is accordingly an object of the present invention to provide an improved printhead having the capability of operating at higher frequencies than in the prior art.\nIt is a further object of the invention to provide a printhead in which the valve/output tube/nozzle has a natural frequency above and sufficiently removed from a desired operating frequency range that there is no resonant interference and consequent degradation in print quality.\nA further object of the invention is to provide a method of optimizing the valve/output tube/nozzle for a printhead of the type described to maximize print quality by adjusting the fluid resonant frequency thereof.\nThese and other objects of the invention will be apparent from the remaining portion of the specification."}
{"text": "The present invention relates to the field of flow guiders or spoilers for wipers and wiper assemblies, and in particular to a flow guider for a wiper and a wiper assembly which are provided with good performances on stability and disassembly to assist the wiper spring rail to be connected firmly to the wiper.\nAutomobiles are currently universal transport used by the public, especially when it is windy and rainy. In addition to act as a replacer of walk, automobiles have a better function of sheltering from wind and rain than bicycles or motorcycles. Automobiles will mostly use wiper to keep driving sight and ensure safety when running in the rain. Generally speaking, the wiper comprises a wiper rod, of which a connecting portion cooperates with an adapter so that a wiper arm is connected to the wiper rod. The wiper arm can be driven by a driving component so as to drive the wiper rod connected to the wiper arm, thus enabling the wiper to wipe the windshield. In other words, vehicle wiper device is typically mounted adjacent to vehicle windshield. The wiper rod of the wiper reciprocates in the range of the wiper of the windshield, whereas the wiper rod is used to remove depositions such as dust or water stains, for example, rain or snow pieces adhered to the windshield.\nThe end face of one end of a portion of existing wiper flow guiders has mostly open configuration. The above-mentioned open end of the flow guiders is used for receiving spring rail, whereas after the spring rail has been received in the flow guider from the above-mentioned open end of the flow guider, it needs to additionally use an end cap to be fixed to one end of the flow guider to enclose the end of the spring rail appropriately, thereby preventing the spring rail passing through the flow guider from the end of the open flow guider. However, such a configuration undoubtedly adds additional manufacturing process and material cost.\nThe other portion of existing wiper flow guiders are configured that the end face of one end of the flow guiders is close in order to improve manufacturing process and material cost problems caused by the additional use of end cap. However, such a flow guider merely close its end face and fails to enclose the end of the spring rail appropriately at the end. It is therefore unable to produce an effective fixing effect to the end of the spring rail. That is, two ends of the spring rail are still easy to break away from the fixing mechanism of the flow guider during using the wiper, resulting in two ends of the spring rail protruding out of the fixing mechanism of the flow guider and the wiper fails to work. In addition, such a wiper flow guider is further provided with a stop portion in the shape of a rectangle, and the end of the stop portion connected to the side surface of the flow guider is the end that has smaller width, thereby to cooperate with the spring rail having a groove in the shape of a trapezoid so as to prevent the displacement between the spring rail and the flow guider in the longitudinal direction. However, the stop portion with such a configuration is not easy to be embedded into the groove, and therefore causing difficulties in assembly. Furthermore, when the stop portion with such a configuration is embedded into the groove, the stop portion will produce a considerable distortion and therefore, the stop portion is prone to be damaged and unable to server the fixing function after repeated assemblies and disassemblies. Therefore, there should be more sophisticated needs for existing flow guiders in the respects of fixing the spring rail firmly and preventing the breakaway of the spring rail and facilitating the user's operation."}
{"text": "Air purification devices are commonly employed in home, office and industrial settings as a way to remove harmful components from the ambient air and to improve air quality. Potentially harmful components that can be found in the air include chemical pollutants, odors, airborne particles, bacteria and/or the like. Air purification devices can be employed for many uses including treating the exhaust gases from motor vehicles and other processes, for the generation of sterile environments and “clean” manufacturing facilities, and for removing pollutants that are commonly found in home, office and industrial environments.\nVolatile organic compounds (VOCs) are gaseous phase organic compounds that are present in the environment from a variety of sources. VOCs are can be outgassed by common manufactured items found in home and office settings such as carpets, paint, cleaners, furniture, and plastics as well as sources such as cigarette smoke or exhaust from passing automobiles. VOCs are also produced across a range of industrial processes including refineries, semi-conductor manufacturing plants, and chemical manufacturing including paints, coatings, pharmaceuticals.\nThe levels of pollutants such as VOCs can be 2-5 times higher indoors compared to outdoor levels. Mechanical ventilation systems can help decrease the levels of indoor pollutants depending upon the levels of pollutants found in the outside environment. A mechanical ventilation system can turnover an indoor air volume as the rate of 1 to 3 hr−1, while the turnover rate can be as low as 0.1 to 0.4 hr−1 without mechanical ventilation. The largest reduction of airborne pollutants including VOCs can be achieved by directly treating air to remove pollutants or to convert pollutants to harmless byproducts."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of optics. More particularly, the present invention relates to an optical system that transforms a light beam having any axially-symmetric intensity distribution to a light beam having another axially-symmetric intensity distribution.\n2. Description of the Related Art\nLasers emitting collimated beams of coherent light have many applications in optical science and technology, including lithography, spectroscopy, communications and display technology. Due to fundamental properties of light propagation in optical resonators, most lasers emit beams having a light intensity that is extremely inhomogeneous. Specifically, the light intensity frequently follows a Gaussian distribution, such as EQU I(r)=(2P/.pi.w.sup.2)e.sup.-2r.sup..sup.2 .sup./w.sup..sup.2 , (1)\nwhere I(r) denotes the optical power per unit area measured at a distance r from the axis of the beam, P denotes the total power of the beam, and w is the beam waist parameter, which sets the length scale over which the optical intensity declines from its maximum value to zero. The same distribution also describes, to a good approximation, the intensity profile of a beam that emerges from a single-mode optical fiber, such as is used extensively in the optical industry for conveying coherent light.\nFor many applications, it is desirable that some area of interest be illuminated as uniformly as possible. For example, optical lithography, which is used to fabricate microelectronic devices, requires that the light fluence over an entire exposed region conform to tight tolerances. Laser users, therefore, frequently encounter the problem of transforming a beam having a Gaussian intensity profile to a so-called flat-top profile, which has a uniform intensity over a region of arbitrary radius a.\nMany solutions have been proposed for transforming a Gaussian beam to a flat-top beam. All conventional solutions, however, have significant drawbacks. For example, the conceptually-simplest conventional method uses an element having radially-varying absorption for removing excess intensity from the center of a beam. Such an approach is inherently inefficient because it can be shown that, in the best case, the fraction of the incident beam power that emerges in the apodized beam is 1/e, or approximately 37%. Moreover, in this conventional approach, the absorptive element only subtends the central part of the incoming beam, having an aperture or other discontinuity located at a point where the light intensity is an appreciable fraction of the peak intensity.\nWhen using spatially coherent light sources, including most lasers, any aperture that truncates the beam also diffracts light into the central region. Accordingly, interference of the diffracted and transmitted light reduces the uniformity of the beam. Yet another drawback of this conventional approach is that stable, well-characterized absorptive materials are required, which are not available for the technologically-important ultraviolet wavelengths.\nAnother conventional approach uses lithographically or holographically fabricated phase gratings for reshaping a Gaussian beam by diffraction. Holographic gratings suffer from limited diffraction efficiency of only about 30%, as well as a lack of materials that are suitable for ultraviolet applications. Lithographically-fabricated phase gratings can have high efficiency, but are expensive to fabricate and only work as designed for a single wavelength. Additionally, it is exceedingly difficult to avoid diffraction into unwanted orders, leading to undesirable effects, such as non-uniformity of the output beam at high spatial frequencies and \"hot spots\" on the beam axis.\nConventional refractive solutions have been proposed that use either spherical or aspheric optical elements for aberrating and then recollimating a laser beam. The solutions with conventional spherical optics are physically bulky and relatively inefficient because the spherical surfaces introduce limited aberrations. More compact and efficient conventional designs have been proposed that use either aspheric or gradient-index lenses. Nevertheless, use of a gradient index accomplishes essentially the same result as an aspheric surface, but having a drawback that no gradient-index glasses are available for ultraviolet applications. All aspheric and gradient-index solutions that until now have been analyzed in detail have used a negative first element and a positive second element in a configuration resembling a Galilean telescope. Unfortunately, such designs require lenses having large deviations from sphericity and are difficult to fabricate. Fabrication problems are especially acute for the concave surface of the first lens in the Galilean design.\nAnother serious problem with most conventional aspheric and gradient-index solutions is that such solutions are only valid for the central region of an incident beam, thereby entailing an aperture for the other discontinuity at a point where there is appreciable input beam intensity, at least 1/e.sup.2 times the peak intensity. As previously mentioned, truncating an input beam causes diffraction and interference fringes that reduce the uniformity of the output beam.\nWhat is needed is a way to efficiently transform a light beam having any axially-symmetric intensity distribution to a light beam having another axially-symmetric intensity distribution."}
{"text": "The present invention relates in general to time-division multiple-access communications and in particular to frequency and timing estimation for burst acquisition of random access channels.\nSeveral communication techniques are well suited for applications where a number of independent users need to share a common band of frequencies without the benefit of an external synchronizing mechanism, such as in time-division multiple access (TDMA) cellular radio and satellite systems. TDMA is a digital radio frequency signaling technique, in which the transmission of a message occupies the communication channel for only a fraction of the total time (i.e., a time slot 22), in a periodic fashion, thus allowing other users access to the same spectrum on a time-shared basis.\nIn order to achieve this sharing of the spectrum, TDMA systems rely on tight synchronization of time bases, between transmitter and receiver, to enable communication. To synchronize the stations, each time slot 22 begins with a reference burst.\nIn an ideal system, the receiver knows that a burst always arrives at the beginning of a slot at a known radio frequency (RF) and prepares to detect it at the beginning of each slot. However, due to a multitude of factors such as propagation delay differences among users at different locations, Doppler shift, and oscillator drift, the offset of the arrival time relative to the expected time and the actual radio frequency of the random access channel burst have a wide range of uncertainty. Therefore, a large detection window in both time and frequency must be processed to look for the burst. Current transmissions begin each data frame (i.e., a complete period comprising one time slot from each transmitter) with a preamble containing a synchronization pattern. The receiving system attempts to detect this pattern, by looking for spectral energy in certain frequency bins (i e., a series of adjacent frequency ranges), and then estimates the frequency and timing using interpolation.\nHowever, these systems are often inaccurate, and extraneous signals (e.g., noise) can cause the receiver to fail to detect the random access channel burst when present, erroneously detect a burst when none is present, or successfully detect the burst but fail to estimate the correct timing or frequency offset from the expected timing and/or frequency offset, thereby losing data. What is needed, therefore, is an improved method of estimating the frequency and timing offsets of a random access reference burst.\nThe present invention relates to communication systems that use a reference burst to synchronize a receiver with a transmitter. The present invention describes a message format and associated method to improve the accuracy of estimating the frequency and timing offset, during initial burst acquisition of a random access channel, in a time-division multiple access communications system.\nIn accordance with one aspect of the present invention, a novel signal structure is shown in which a continuous wave (CW) segment used primarily for frequency estimation is separated into at least two parts separated by some portion of the content of the random access channel burst to provide more precise frequency estimation. Increased precision in frequency estimation is achieved as a result of having longer periods of time over which to measure the time rate of change of the phase of the CW signal.\nIn accordance with another aspect of the present invention, iteration of separate frequency and timing estimation procedures is used to refine both estimates, as the individual estimation processes are more accurate when there is less error in the other estimate. This method detects the presence of random access channel bursts by correlating the frequency of the continuous wave segment of an incoming data packet with the expected frequency of the burst. This is accomplished with a plurality of filters near the expected frequency, configured to alert the system when a preset threshold is exceeded. The center frequency of the detection filter provides the receiver with a coarse estimate of the unknown frequency offset. A coarse timing estimate is obtained by continuing to monitor the detection filter after the alert and determining the sample point of the peak filter response. The frequency estimation is refined by interpolating the filter response from the detection filter and its neighbors at the acquired sampling time. The timing estimation is refined by correlating the frequency corrected received data against a template representing both the continuous wave and unique word segments. Because the accuracy of the fine timing estimation depends on the residual frequency error and vice versa, the fine timing and fine frequency estimations are performed iteratively.\nIn accordance with yet another aspect of the present invention, multiple hypothesis testing is used, in which more than one initial frequency estimate is carried through the iteration process, with thresholding to identify the best frequency offset. Reliable rejection of false ambiguities is accommodated since the degradation in timing estimation due to large frequency offsets can be easily detected with thresholding.\nThe present invention provides significant increases in the precision of timing and frequency offset estimations, when acquiring random access channels using a reference burst, by separating the continuous wave portion of the message, iterating between frequency and timing estimates, and carrying multiple hypotheses through the iteration process."}
{"text": "The present invention relates generally to bearing cages for use in rotating machinery and, more particularly, to bearing cages molded of a polymer.\nBearing cages are conventionally used in ball and roller bearings to retain and separate rolling elements (rollers, needles or balls) in predetermined relative positions. The bearing cages lie between inner and outer bearing races.\nIn certain applications, bearing cages are exposed to relatively difficult conditions that limit the useful life of the bearing cages or require the use of complex and expensive construction. One such application is in the needle bearing assembly for a crank of a two-cycle engine such as, for example, an outboard motor.\nIn such an application, the needle bearing assembly is exposed to high loads generated by combustion pressure, inertia of the parts, and centrifugal force. This application is so severe that a special high precision bearing cage made of silver-plated steel, is generally used. The silver is believed to act as a lubricant, at least during a break-in period of the bearing.\nThere are numerous disadvantages to the above-described bearing cage. The high precision requires numerous steps for manufacture, such as piercing, heat treating, grinding, and chemical finishing. The initial cost of tooling is very high, and the chemical by-products produced during manufacture result in environmental problems of disposal. In addition, because of the many operations which must be performed, the time for manufacture of such bearing cages is long. Thus, a larger than desired inventory of work in process must be maintained. The large number of operations, combined with the required precision, may result in high labor costs or a relatively high scrap rate. And, the required silver substantially adds to the cost.\nThe foregoing illustrates limitations known to exit in present bearing cages. Thus, it is apparent that it would be advantageous to provide an alternative directed to overcoming one or more of the limitations set forth above. Accordingly, a suitable alternative is provided including features more fully disclosed hereinafter."}
{"text": "This disclosure relates to a cushion structure, assembly and its construction for comfort applications, including seating and bedding utilizing polyester fiber as one filler material.\nA typical cushion assembly used in seating applications includes an aesthetic cover surrounding a soft and resilient filler material such as polyurethane foam, springs and the like. Most cushions are constructed with material which provides a desired support and comfort to the user. Polyester fiber toppers are sometimes used on top and side of the cushion assembly on which a user may sit to provide a better “hand,” which is a desired feel. An additional wrap is occasionally used to provide a desired function such as to provide improved resistance against flammability. The wrapped and/or padded core structure, usually made up of polyurethane foam is inserted into an aesthetic cover. The foam core is generally the same dimensions as the cover or very slightly larger than the cover.\nThe feel of foam cushions is very customizable. This is done by changing the foam chemistry for a given density. One measurement of “feel” for a cushion is the Indentation Load Deflection, ILD, which is determined using industry guidelines. The ILD is the amount of pounds (measured as resistant force) required to compress a 4 inch thick, 15 inch×15 inch sample to 3 inches (or 25% of original height). For example, a typical 4 inch tall polyurethane foam cushion having a density (in pounds per cubic foot, or “pcf”) of 1.0 pcf has an ILD of 30. In addition due to processing and chemistry changes this is tunable within a range of 10-40 ILD; a density of 1.2 pcf is tunable to 20-50 ILD; and so on.\nA given foam cushion must also exhibit an acceptable comfort or “support factor,” typically in the range of 1.7-3.0. The support factor calculated by dividing the force required to compress a 4″ sample to 65% of its height by the force required to compress to 25% of its height; i.e. comfort factor=(ILD @ 65%)/(ILD @ 25%) In addition to comfort, a standard foam cushion must survive industry durability tests over several thousand cycles during which the foam cushion must substantially maintain its height and shape while maintaining the support factor. Polyurethane foam cushions are highly tunable in that a foam material can be easily selected to provide a desired density, ILD and support factor, which in turn provides the durability for a given application. For polyester fiber cushions, the ILD is very closely tied to the density and such fiber cushions are not easily tuned to provide both comfort and durability. To provide better durability fiber cushions must be made very dense but that is generally not acceptable as comfortable for the end user.\nFoam is generally very resilient and tunable but the chemistry is such that in its native state it is very highly flammable, further the process of making foam is considered harmful to the environment. Several attempts to replace foam with polyester fiber have resulted in different formed, thermo-bonded or loose fiber constructions, but none have been able to achieve desired comfort and performance sufficient for broad commercial viability due to the difficulty in tuning. Loose fibers or thermo-bonded fibers are sometimes used in outdoor cushion applications to provide improved long-term resiliency over foam. The fibers are loosely arranged relative to one another in an unconnected fashion, permitting the fibers to shift uninhibited within the cover. Such outdoor cushions are not very durable and as such not suitable for conventional seating and bedding applications.\nTufting has been used to secure multiple layers to one another or provide an aesthetically pleasing exterior cover. Often, the tufts are visible through the cover. For example, tufting is used in futons to secure the multiple layers to one another and the exterior cover. The exterior cover is arranged rather loosely about the layers."}
{"text": "The function detection of a loudspeaker module not only involves perfection detection of a single interface, but more importantly, whether the function of an audio module in loudspeaker front-end device is normal is judged by detecting the signal of a loudspeaker interface, wherein the audio module includes an audio generating source module, an audio transmission channel and an audio control module. Thus, the function detection of the loudspeaker interface is substantially to judge whether the function of the audio module of the loudspeaker module is normal. Before overall assembly, the function of the loudspeaker module of a product in which the loudspeaker is arranged is detected, so that defective audio module products may be effectively prevented from entering the subsequent assembly stage, then the reject ratio of the overall product is controlled, and the production cost is reduced.\nAt present, the most common method for detecting the function of the loudspeaker module includes the steps of directly connecting the loudspeaker to a detection interface and then judging whether the loudspeaker produces sound by human ears to determine whether the function of the audio module is normal. This method is most direct and convenient. However, generally, a plurality of loudspeaker modules need to be simultaneously detected in the detection process. Accordingly, the sound simultaneously produced by the plurality of loudspeaker modules is liable to generate acoustic interference to human ears, so that the detection result is inaccurate and the detection process needs a large amount of manpower."}
{"text": "1. Field of the invention\nThis invention relates to polymeric dyes obtained by condensation of cationic biopolymers with dyes, to a process for their production, to formulations containing these dyes for dyeing fibers, to processes for dyeing fibers using the polymeric dyes and to their use for dyeing fibers.\n2. Discussion of the Related Art\nSo-called substantive dyes or oxidation dyes are normally used for coloring keratin fibers, preferably human hair. Oxidation dyes consist of a primary intermediate and a secondary intermediate and are formed on the hair itself under the effect of oxidizing agents or atmospheric oxygen.\nHowever, there is a need on the market for coloring products which enable the color of hair to be rapidly changed, i.e. which can be quickly and completely washed out. The hair normally has to be washed 5 to 6 times to remove known products, something which many consumers find unsatisfactory.\nAnother disadvantage of known hair coloring or hair tinting formulations is that the products always penetrate into the keratin fibers and, hence, cause damage--albeit normally slight--to the hair structure. Since the formulations do not act selectively on keratin fibers, the treated scalp is generally colored at the same time which consumers also find undesirable.\nAccordingly, the problem addressed by the present invention was to provide new formulations for the temporary coloring of keratin fibers which would not have any of the disadvantages mentioned above."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for signal transmission/reception in a communication system using a Hybrid Automatic Repeat reQuest (HARQ) scheme.\n2. Description of the Related Art\nNext generation communication systems are developing towards systems capable of providing a high speed large-capacity service to Mobile Stations (MSs). Representative examples of the next generation communication systems include an Institute of Electrical and Electronics Engineers (IEEE) 802.16 communication system, a Mobile Worldwide Interoperability for Microwave Access (WiMAX) communication system, and an IEEE 802.11 communication system. The Mobile WiMAX communication system is a communication system based on the IEEE 802.16 communication system, and the IEEE 802.16 communication system is a communication system using an IEEE 802.16 standard.\nIn the next generation communication systems, use of Low Density Parity Check (LDPC) codes is taken into consideration to support high speed large-capacity data transmission/reception, because it is possible to achieve parallel implementation of the LDPC codes, which makes it possible to implement a decoder having a throughput higher than other channels codes, such as turbo codes or convolutional codes.\nHereinafter, a conventional parity check matrix for structured LDPC codes will be described with reference to FIG. 1.\nFIG. 1 illustrates a conventional parity check matrix for structured LDPC codes.\nReferring to FIG. 1, structured LDPC codes may be generated using a parity check matrix having a size of mL×nL. That is, each row of the parity check matrix for structured LDPC codes includes m square matrixes each having a size of  and each column of the parity check matrix for structured LDPC codes includes n square matrixes each having a size of  wherein the square matrix may be a zero matrix or a permutation matrix. As used herein, each row or each column of the permutation matrix includes a weight of one (1), wherein the weight corresponds to the number of elements having a value (e.g. 1) other than zero (0). In FIG. 1, Pm,n refers to a square matrix located at a position at which the mth square matrix row and the nth square matrix column intersect when the parity check matrix includes m square matrix rows and n square matrix columns.\nAccording to the parity check matrix for structured LDPC codes as described above with reference to FIG. 1, it is possible to generate LDPC codes having various lengths by using only one parity check matrix, through a change of the size (i.e. the vale ) of each square matrix in the parity check matrix. Hereinafter, the scheme of changing the size of the square matrix to generate LDPC codes having different lengths will be referred to as a “lifting scheme”.\nMeanwhile, in the next generation communication systems, various schemes including a Hybrid Automatic Repeat reQuest (HARQ) scheme and an Adaptive Modulation and Coding (AMC) scheme have been proposed to support high speed large capacity data transmission/reception. To use the HARQ scheme and the AMC scheme, it may be necessary to support various code rates.\nHowever, the structured LDPC codes are disadvantageous in view of the code rates. That is, since the structured LDPC codes are generated by using a parity check matrix having a fixed size, it is relatively difficult to change the length of the structured LDPC codes, in comparison with other channel codes. Therefore, the structured LDPC codes are not flexible in view of the coding rate."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus for evaluating the mechanical properties of materials that can evaluate the mechanical properties of a test specimen in magnetic fields of two different orientations.\n2. Description of the Prior Art\nIn superconducting generators, superconducting magnetic energy storage devices, linear motors and other such superconducting device technologies, the importance of the relationship between the mechanical properties and the magnetic field orientation of superconducting materials and structural materials is such that prior evaluation of this relationship is required. However, the difficulty of making an apparatus for evaluating data relating to the mechanical properties of such materials has meant that the amount of such data that has been collected is insufficient.\nIn a conventional apparatus for evaluating the mechanical properties of a material in a magnetic field, a test specimen is placed in a superconducting magnetic field and subjected to a load while the strength of the material under the load is evaluated. Such apparatuses for evaluating the mechanical properties of a material having a unidirectional magnetic field orientation have already been announced. For example, on page 119 of the proceedings of the 45th meeting on Cryogenics and Superconductivity of the Cryogenic Association held in the spring of 1991, a system arrangement is shown in which a solenoid type magnet in which a relatively high magnetic field can be readily generated is used to apply a load in order to evaluate the properties of a material in a magnetic field oriented in the same direction as that of the direction in which the load is applied.\nHowever, while a solenoid type magnet can be used to evaluate the mechanical properties of a test specimen in a unidirectional magnetic field, the magnetic field orientation cannot be changed. This makes it impossible to evaluate the mechanical properties of a material under different magnetic field orientations, and as such makes it inadequate for evaluating the mechanical properties of materials that are to be used in fields relating to superconduction technology. While for the purpose of evaluating the mechanical properties of a material in different magnetic field orientations a separate superconducting magnet may be provided to generate a magnetic field in a different direction, such an arrangement would require the provision of at least two superconducting magnets, which would thereby increase the size and cost of such an apparatus."}
{"text": "The magnetic sensor can be used to sense the magnetic components of an external magnetic field such as the geomagnetic field. The magnetic sensor has been applied to fields such as automobiles, auto-control systems, medical systems, and electronic compasses, in which the electronic compass needs to sense the three magnetic components of the external magnetic field. The smartphone is typically implemented with an electronic compass nowadays. The electronic compass needs to sense the geomagnetic field and thus requires the magnetic sensing capability to sense the three axes of the magnetic field. In addition, high sensitivity is needed for the electronic compass because of the geomagnetic field.\nThe magnetic sensor may have various designs. The magnetic tunnel junction (MTJ) device has been widely used to sense the horizontal component of an external magnetic field. The magnetic sensing layer of the conventional MTJ device may be formed by in-plane magnetic anisotropy (IMA) materials, which have relatively high sensitivity for the in-plane component of the magnetic field but relatively low sensitivity for the out-of-plane component of the magnetic field. In other words, even if the thickness of the magnetic sensing layer for the conventional MTJ device is controlled at the superparamagnetic condition, the MTJ device is still insensitive to the out-of-plane (Z-axis) component of magnetic field with relatively high sensitivity for the in-plane (X-axis and Y-axis) components of the magnetic field. One reason is that, although the magnetic moments for the superparamagnetic IMA are isotropically distributed, the superparamagnetic IMA maintains the in-plane property. Therefore, the response to the out-of-plane (Z-axis) magnetic field needs to overcome the demagnetization energy in order for the magnetic moments to be perpendicular to the magnetic sensing layer, resulting in insensitivity to the Z-axis magnetic field.\nA MTJ device with a simple structure for sensing the out-of-plane (Z-axis) magnetic component of the geomagnetic field is under development."}
{"text": "This application claims the priority of German patent Application No. 199 16 452.5 filed Apr. 12, 1999, which is incorporated herein by reference.\nThis invention relates to a power ring that has an annular conductor structure and is installed in a vehicle, particularly a military vehicle, for supplying devices in the vehicle with electric energy. The power ring includes a plurality of controllers functioning as monitoring and switching devices that are connected along and on an annular conduit, as well as a plurality of taps on the annular conduit for connecting the respective devices with the annular conduit. The power ring includes feed locations along the annular conduit for supplying electric energy and data connections between the controls, and a superordinated control device for monitoring and controlling the power ring.\nReliably supplying electric energy to devices which are individual components that consume electric power in a vehicle, particularly in a military track laying or wheeled vehicle, is effected increasingly by means of a xe2x80x9cpower ringxe2x80x9d in which, in contrast to a star-like connection of each load (electric energy consumer) to a generator, an annular conduit that conveys the energy connects each electric load with the current producer or the generator.\nA short circuit on the annular conduit of the power ring results in an interruption of energy supply of all the loads connected to the power ring. Then, conventionally, a repair workshop is visited to eliminate the defect. In particular for applications in military vehicles, however, an interruption of the internal current supply to all devices connected to the power ring because of a simple short circuit is undesirable and not permissible.\nConventional solutions provide to subdivide the power ring into individual segments by means of a power ring controller. The devices are connected to the segments of the power ring by taps on the power ring. In case of a short circuit in one segment, only those devices that are connected to the respective segment are affected by the short circuit and become inoperative. The devices connected to the other segments of the power ring continue to be supplied with electric energy and thus continue to remain functional.\nGerman patent document 197 33 760 discloses, in connection with a data bus for information transmission, that transmitter and receiver components connected to the bus conduit at nodes may switch on a terminal resistance at each node, and further switches may be placed at the node so that a device connected to the node may be disconnected from the bus conductor in case of a defect, while the remaining bus operation is not affected by the failure of one node.\nThe conventional solutions have the disadvantage that by dividing the power ring into ring segments, in case of a failure the defective power ring segment is switched off by means of the power ring controller, as a result of which the devices connected to such a segment are no longer supplied with electric energy. In case of a small number of segments as well as controllers and the same number of devices, such a disadvantage is even greater because a relatively large number of devices fail with a single segment because a greater number of devices per segment is supplied with electric energy.\nA redundancy construction of the power ring may be remedied for a device breakdown in case of individual failures. Such a construction, however, requires substantial outlay and is disturbance prone because an additional failure sensing has to be provided and the number of the power ring components are doubled.\nIt is an object of the invention to limit the effect of a short circuit in a power ring in such a manner that only a few connected devices fail, that is, a possibly large number devices are not affected by a simple short circuit. A simple possibility should be found to continue to supply all devices with electric energy despite a breakdown of a power ring segment also viewed as a connecting conductor, at least for failure number one (single failure case).\nThe above object generally is achieved according to the present invention by a power ring of annular conductor configuration for use in a vehicle, preferably a military vehicle, to supply electricity to devices in the vehicle, with the power ring including: a plurality of controllers as monitoring and switching devices arranged along and connected with a plurality of conductor segments connected to form an annular conductor; a plurality of taps provided for connecting the annular conductor to loads; at least one feed location connected to the annular conductor for feeding electrical energy to the annular conductor and to the loads; and data connections from the controllers to a superordinated control apparatus for monitoring and controlling the power ring; and wherein: each tap and feed location for the annular conductor is structurally combined as a unit with a controlled node of a respective controller; each tap controller is connected along the annular conductor and into the annular conductor in a multiple manner so that the annular conductor segments are disposed between adjacent tap controllers, and at least one load or a current feed connected to a tap controller is selectively connectable with adjacent conductor segments of the annular conductor via a plurality of controlled switches disposed at the tap controller of the respective controller with which the respective annular conductor segments may selectively be connected to or disconnected, and with which the respective load or the feed may selectively be connected or disconnected.\nAccording to the invention, the power ring is provided with its structural elements in such a manner that the power ring controller and the power ring tap are combined into a common system unit and device. The thus-formed tap controllers are interconnected with conductors as a power ring in such a manner that the devices to be supplied with electric energy need only to be connected to the tap controller. In case of a failure, one segment of the power ring which is situated between two tap controllers, is switched off, and all devices which are connected to the tap controllers may be continued to be supplied. In case the failure concerns only the power ring conductor, all devices continue to be supplied because these are only connected to the tap controllers and the annular conductor continues to be closed via the remainder of the ring. Only in case of a failure of a tap controller is that portion of the devices which is coupled to the controller no longer supplied. Such a condition would appear also in case of two conductor segment failures. In case of a single conductor segment failure or isolated individual failure, however, all devices continue to be supplied by the power ring according to the invention."}
{"text": "Staircases are often used to allow passengers to walk from the tarmac to the passenger compartment of an aircraft. Similarly, staircases allow cargo handlers to walk from the tarmac into a cargo hold of the aircraft. The staircases used to access the passenger compartment or cargo hold of the aircraft are typically stored in a hangar or on the tarmac at an airport facility. The staircases are also often equipped with wheels so that they may be transported from their storage locations to adjacent the aircraft. In some instances the staircases may be mounted on a vehicle capable of travelling from a hangar or other waiting area of an airport to the aircraft. The vehicle may also be capable of raising the staircase to a height suitable for a person to access the aircraft passenger compartment or cargo hold via the staircase.\nSuch movable staircases are heavy, expensive, and also require at least a paved or unpaved pathway over which the staircase can be moved from its storage/parking area to the aircraft. Although this is not generally a problem at many modern airports, some airports in remote and/or undeveloped parts of the world may still not be equipped with wheeled staircases and/or with adequate pathways that may allow a wheeled staircase to be transported from a hangar/storage area to the aircraft.\nOne example of an aircraft that may travel to such remote and underdeveloped airports is the Flying Eye Hospital of Orbis (http://www.orbis.org). The flying eye hospital is housed in an aircraft, in which the passenger compartment serves as a mobile eye clinic, and the cargo compartment serves as a storage area for medical supplies and other equipment. Once the Orbis aircraft lands at an airport, patients will require a staircase and/or ramp to climb up from the tarmac to the passenger compartment to receive treatment. Likewise, medical personnel will require a staircase to deplane and/or climb up from the tarmac to the cargo hold to retrieve medical supplier or equipment.\nThere is a need for a staircase that can be easily disassembled and stored in the cargo hold of the aircraft, particularly when the aircraft travels to remote or underdeveloped airports or landing areas. This may allow the staircase to be used even at airfields that may not be equipped with wheeled staircases or may not have ground facilities to provide staircases to access the passenger or cargo compartments of an aircraft. The modular staircase of the present disclosure represents an improvement over the prior art and solves one or more of the problems set forth above and/or other problems of the prior art."}
{"text": "In new consumer electronics devices, there are increasingly larger arrays of sensors for monitoring the performance and functionality of the device and making diagnostic decisions. Automobiles are perhaps the most advanced in this area. Though not traditionally retail consumer items, aircraft and elevators and escalators have been developed in the past decade with a substantial number of sensors built into them so as to better diagnose problems in their operation. The additional cost of manufacturing due to the added sensors is generally seen by the manufacturers as more than offset by more reliable performance (the sensors aid in anticipating problems before they become large enough to cause major failure), decreased down time due to faster and more accurate diagnostics, and a more competitive position for the original manufacturer to secure long term maintenance contracts for the equipment it originally sold.\nHowever, this trend toward sensorification is not proceeding as quickly in lesser-value manufactured goods, and the vast majority of other consumer electronics with engines still do not have very many diagnostic sensors if any at all. Much like auto repair was routinely done a generation or two ago in what are now more developed countries, today there is a lack of highly specialized diagnostic equipment for auto repair in lesser developed countries. Sometimes even rudimentary electronic diagnostics like an engine analyzer or a timing light are unavailable. Cars being repaired are older on average than those being serviced in more developed countries, and even if they did include the inbuilt diagnostic sensors noted above there is no electronic equipment to extract and exploit that information. This less refined diagnostic process is likely to persist for some years even in the face of favorable conditions of economic growth across the local populace. The problem is not limited only to the automobile repair industry though.\nFurther, in more developed countries there is a concerted effort to make various devices electromagnetically compatible, such that one device does not cause electromagnetic interference EMI with other common devices. This can be seen readily in portable or auto-mounted radios. In the past EMI between an AM radio and a cellular phone was an issue, whereas largely that problem has been solved and there is no need to turn off an AM radio when communicating on a cellular phone. In less developed countries EMI has not been an historic problem and so devices there are not prone to such EMI compatibility as they are in areas where electronics have become more ubiquitous.\nWhat is needed in the art is a way to remotely diagnose products in less developed areas of the world, as well as lower-value consumer devices that do not have in-built sensors similar to those currently used in high-cost manufactured items such as automobiles, aircraft and elevators/escalators. It would be advantageous to address this diagnostic problem in a manner that does not require a high capital investment in diagnostic equipment. The lower cost and relatively lower complexity of the vast majority of electronic devices that are directed to the retail purchaser/end user generally does not justify having such an array of sensors built into the manufactured device, yet the complexities are high enough that diagnosing a problem can be time consuming and expensive. For this reason especially in the developed world many consumer electronics are deemed more ‘disposable’ than need be, because often the cost of diagnosis is high enough to make diagnosis and repair an uneconomical option. An economical solution to the diagnostic problem would change the balance of the questions whether to repair or replace a malfunctioning device."}
{"text": "1. Field of the Invention\nThe present invention relates to a variable speed power tool, more particularly to a power tool (such as an electric drill, a screwdriver or an impact drill) having a control system for automatically varying the output speed in accordance with a load applied to the power tool.\n2. Description of the Related Art\nA variable speed power tool (such as for example an electric drill, a screwdriver or an impact drill) usually employs a gear transmission system to selectively output various speeds to adapt to different work conditions. The gear transmission system generally includes a gear transmission mechanism connected between an electric motor and an output shaft and a control system operable by the user to change the gear transmission mechanism so as to vary the gear reduction ratio to output various speeds.\nU.S. Pat. No. 6,431,289 discloses a gear transmission mechanism including two ring gears being adjustable forwards and backwards via an actuating member operated by the user outside the housing. By this means, the ring gears can shift between a rotating position where they connect both planet gears and an adjacent planet carrier and a fixed position where they connect to the planet gears and the housing thereby varying the output speed. However by such manual operation, the operator has to monitor the running bit (such as a fastening screw) to decide when to vary the output speed or torque to achieve better output efficiency. This operation is not user-friendly. Furthermore, when the user actuates the actuating member to shift the ring gear from the rotating position to the fixed position, the rotating ring gear is unable to engage a fixed structure in the housing smoothly. In these circumstances, a gear clash between the ring gear and the rotationally fixed structure frequently occurs which results in damage to the ring gear or the fixed structure and may reduce its working life.\nU.S. Pat. No. 6,824,491 discloses an automatic variable speed system in which a control system is able to automatically vary the gear reduction ratio by varying the direction of movement of the ring gear according to variation of the torque load. EP-A-0787931 discloses a similar automatic variable speed system. However, these automatic variable speed systems are performed mechanically and the mechanical structure of the systems are complex and difficult to manufacture. Neither system is able to eliminate the occurrence of gear clash.\nP-A-3221381 discloses a solenoid mounted on a motor which is able to actuate linear movement of a switching lever via an actuation lever and a transmitting lever to change the gear reduction ratio of the gear transmission mechanism.\nJP-A-4262151 discloses a solenoid adapted to change the gear reduction ratio of the gear transmission mechanism.\nJP-A-8068461 discloses a keep solenoid used to drive an internal gear from a high speed position to a low speed position and retain the internal gear at the low speed position. A spring is used to pull the internal gear back to the high speed position.\nP-A-9057639 discloses an auto speed change mode which will not be activated if the motor does not reach the highest output. It is undesirable that auto speed changing is actuated when the trigger switch is being depressed to increase the motor speed. A motor speed detect module is provided to detect motor speed and feed a signal representing motor speed to the control circuit.\nJP-A-9057640 discloses a mode selection switch provided to activate or inactivate an automatic speed changing function. The switch is used to select the auto speed change mode and the fixed speed mode. The switch can select between auto mode, high speed mode and low speed mode.\nJP-A-10103462 discloses a curve ηh indicating the tool output efficiency at a high speed state. ηl is a curve indicating the tool output efficiency at a low speed state. The tool first works in the high speed state and when the load torque reaches CP, the auto speed change is activated so that the tool then works in the low speed state.\nJP-A-9014433 discloses a keep solenoid which is used only to change from high speed to low speed and various means for returning to high speed.\nJP-A-1109845 discloses a battery pack voltage detector VP which will detect BP voltage. When the voltage drops below a predetermined value, the auto speed mode is activated directly to vary from high speed to low speed. The activation of auto speed mode is not controlled by control circuit C. The control circuit C detects the load torque by monitoring the motor speed. If the power drops rapidly and the motor speed reduces sharply, the control circuit C cannot activate auto speed changing before the motor has stalled."}
{"text": "Until recent years, most of the mechanisms used in various package and product dispensing machines relied on a multitude of motors, switches and solenoids for moving various machine parts and otherwise handling the package to be dispensed. Most such machines required one motor, switch and/or solenoid for each row, column or type of product or package dedicated to kicking out or otherwise ejecting the desired package or product therefrom. Therefore, a machine having an arrangement of nine racks of product would typically require nine motors, switches and/or solenoids for ejecting one product from each rack.\nSuch machines generally suffer from numerous disadvantages. One disadvantage lies in the fact that any machine having such an arrangement will commonly experience numerous mechanical failures due to such factors as age, exposure, poor maintenance, poor assembly, misalignment, etc. Furthermore, the tendency to malfunction can be greatly increased when the mechanical parts must operate in a frozen environment such as that created in a frozen food or ice cream vending machine.\nAdditionally, each motor, switch and/or solenoid can be expensive and therefore can contribute to increasing the overall cost of a machine, i.e., either in capital or assembly costs. Elimination of any such mechanical parts will usually result in a lowered cost of production and future maintenance and repair.\nMoreover, extraneous machinery uses up valuable space which could be occupied by product. The loss of product storage space means that the machine location must be visited by the operator more frequently for restocking, thereby increasing operating costs and reducing profitability.\nFor the above reasons and others not expressly discussed herein, numerous attempts have been made to develop package and product dispensing machines which have as few motors, switches, solenoids and other mechanical devices as possible.\nAs one example of such a package and product dispensing machine, consider the apparatus disclosed in U.S. Pat. No. 4,787,533 to Haroutel et al. The dispenser disclosed by Haroutel et al. dispenses postal packets of differing sizes and types. The device utilizes individually controlled and motorized horizontal and vertical carriages for positioning a product holding device in font of the desired article. However, this device still has a relatively large number of electro-mechanical parts which are relatively intricate to assemble and expensive to operate and maintain. Furthermore, the disclosed mechanism is not designed for cold temperature operation and therefore is not suitable for cold food dispensing.\nU.S. Pat. No. 5,025,950 to Trouteaud et al. discloses an apparatus for storing and dispensing ice cream packaged in uniform containers which are stored in rotatable transporters or revolving racks. However, this apparatus does not solve the above-mentioned problems. Firstly, each rotatable transporter requires its own motor drive. Secondly, this arrangement has a very limited capacity. Thirdly, and perhaps most importantly, all of the mechanics of the disclosed apparatus are exposed to the cold environment.\nAs was mentioned briefly above, mechanical failures are common with vending machines that refrigerate their products, especially with vending machines that employ motors or solenoids to dispense their products from the refrigerator. Because the motors and solenoids are located within the refrigerator, they are subjected to cold temperatures and condensates. Such harsh operating conditions increase the likelihood that the motors and solenoids will malfunction. The apparatus of Trouteaud et al. is as susceptible as the other prior art machines and furthermore suffers from the other disadvantages discussed hereinabove.\nGerman Application No. 2,455,673 discloses a machine for vending frozen food packages which utilizes a motorized sliding carriage to position a suction device over the package to be dispensed. The entire suction device is then lowered by a motorized drum and cable system until contact is made between the suction device and the package. The suction device is then activated and withdrawn to a dispensing position along with the package. However, this device still suffers from the disadvantage that the entire mechanism is still subject to the cold environment of the freezer, thereby inviting increased mechanical failures as well as having the inefficiency of expending a good deal of energy cooling large volumes of non-storage space."}
{"text": "The new cultivar ‘PPP CONT 09’ is a product of a cross between the male parent ‘GEDI ONE SAP’ (U.S. Plant Pa. No. 13,731) and the female parent ‘GEDI ONE CAR’ (U.S. Plant Pat. No. 13,794). The new cultivar was discovered and selected by Dirk Pieters in October of 2008 in a greenhouse in Oostnieuwkerke, Belgium.\nAsexual reproduction of the new cultivar ‘PPP CONT 09’ by terminal cuttings was performed in Oostnieuwkerke, Belgium and Oxnard, Calif., and has shown that the unique features of this new cultivar are stable and reproduced true to type on successive generations."}
{"text": "During the storage of liquid phases such as oils, crude, refinery products and other petrochemical products, undesired sedimentation and thickening often occurs which will be discussed hereinafter in connection with the particular deposit-forming example of crude oil.\nThe liquid phase of crude oil is a mixture mainly consisting of hydrocarbons such as paraffins, aromatics and naphthenes, which are also accompanied by non-hydrocarbons or so-called impurities such as mud, water, dissolved salts, sulphur compounds, sand, etc. In certain circumstances prior to processing in refineries, the crude undergoes rough refining processes for the separation of impurities. It is then generally customary to store the crude oil to be processed and also the precleaned crude oil in large tanks. This involves holding times of varying lengths, which can be quite long in the case of stockpiling, and much shorter for operation-based storage locations.\nLong holding times especially favor undesired deposit formation from crude in tanks. These deposits can be of different types, being e.g. favored by emulsions of water with hydrocarbon fractions, or can consist of segregates of heavy hydrocarbon fractions (hard waxes) or settled mud or salts. The result is a kind of oily mud which accumulates and is compressed on the bottom of tanks and leads to high costs and losses and which is referred to loosely as sludge.\nCosts and losses are caused because the oil sludge reduces the capacity of the tanks and also binds crude oil or, in some cases, largely consists of thickened crude. Thus, apart from the costly space loss in the tank, storage leads to significant loss of raw material. In addition, the lost space cannot be recovered again if the sludge is stored in closed systems, i.e. tanks made available for this purpose, the undesirable alternative being a final storage which is prejudicial to the environment if the sludge is dumped into open pits or basins. In large tanks with a diameter up to 100 m, a height of 20 m and a corresponding capacity, sediment thicknesses of 1 to 2 m lead to a 5 to 10% capacity loss. In addition, the service and operation of the tanks is often made difficult because the oil sludge clogs the pumps, outflows from tanks have to be filtered, etc. Finally, down-times are linked with the removal of the oil sludge. If the oil sludge is not removed, it accumulates further and finally leads to the abandonment of sludged-up storage containers and the construction of new tanks. Apart from these storage costs the unprocessed oil sludge also represents a loss because, despite its impurities, it largely consists of utilizable hydrocarbons.\nSeveral solutions are known for removing sediments from crude in tanks and two examples will now be given.\n1) A first solution is proposed in U.S. Pat. No. 3,436,263 and French Patent 22 11 546, where cleaning substances are used for dissolving, or removing in bound form, the oil sludge. A disadvantage of this method is that, due to the cleaning substances which have been introduced, the liquified oil sludge is no longer usable because its composition has been changed by the additives and it must be disposed of in dumps or elsewhere. Such dumps are e.g. old tanks or so-called wasteland and constitute a serious pollution of the environment. Reprocessing of oil sludge is consequently not desirable using this method which, instead of combating the problem it only combats the effect. However, it is still possible to clean and reuse the tanks. The essential reason why the dissolved oil sludge cannot be processed is that the cleaning substances used represent impurities for processing in refineries, whose separation by standard cleaning methods involves great effort and is expensive and does not bear a positive relationship to the recovered crude.\n2) Another solution is proposed in European Patent 160,805, wherein hydrokinetic energy is used in order to dissolve, or suspend back into the liquid phase by means of turbulence, sedimented residues in tanks. Thus, oil sludge dissolved by crude as the dissolving substance can be returned to the process and processed in the refinery following standard cleaning procedures. This method does not prevent the formation of oil sludge but instead merely eliminates it. For this purpose, planned turbulence or eddy flows are generated, whose successive remote action is able to dissolve the deposits in an effective manner, even outside the direct injection zone. However, this requires a considerable investment. Mechanically moving components within the tanks, such as e.g. rotary liquefying lances, which hydro-kinetically activate the oil sludge and dissolve it in crude represent considerable expense. Thus, although this method leads to a high oil sludge recovery level, it is expensive. Under extreme environmental conditions, e.g. in sand or ice desert regions, this is undesired."}
{"text": "With the development of society, information exchange becomes increasingly important. To implement timely information exchange, varieties of instant messaging software, such as WeChat and QQ, are developed one after another. The instant messaging software is mainly used depending on the Internet. Therefore, quality of the network directly affects communication quality of the instant messaging software.\nAt present, with the widespread popularization of network devices, load pressure on the network also increases. When a large number of devices use the network at the same time, network congestion is caused, and affects network communication of network devices. For example, instant messaging of instant messaging software is affected, resulting in poor communication quality between users.\nFor the technical problem of poor instant messaging quality caused by network congestion in a related technology, no effective solution has been proposed at present."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a method and an apparatus for manufacturing a semiconductor device, and particularly relates to a method and an apparatus for manufacturing a semiconductor device for implementing a wafer grinding process on a wafer fixed on an adhesive tape.\n2. Description of the Related Art\nA semiconductor manufacturing process includes a wafer process and a subsequent wafer grinding process. In the wafer process, semiconductor elements are formed on one surface (a front surface) of a wafer. In the wafer grinding process, the other surface (a back surface) of the wafer, without any semiconductor element, is ground so as to reduce a thickness of a wafer. However, it is difficult to handle the wafer with reduced thickness since the wafer may easily break.\nAccordingly, a method and an apparatus for manufacturing a semiconductor device is required, which makes it easier to handle the wafer with reduced thickness and can prevent the breakage of the wafer.\nAlso, because the semiconductor manufacturing process is ever complex, more devices are required for manufacturing semiconductors. Therefore, there is a need for improving space efficiency of a semiconductor device manufacturing plant. Thus, further miniaturization of the apparatus for manufacturing semiconductors is also required.\nReferring now to FIGS. 1-6, a method and an apparatus of the related art for manufacturing a semiconductor device will be described. FIG. 1 shows a prior device structure for a back grinding process that performs wafer grinding on the back surface of the wafer. FIGS. 2-6 show a method of implementing a prior back grinding process.\nAs shown in FIG. 1, in the prior back grinding process, the back grinding process is performed on a wafer 14 using a supply cassette 2, a storage cassette 4, a transfer machine 6, a wafer grinder 8, an ultraviolet (UV) irradiation device 10 and a remover 12. Hereinafter, each of the above devices will be described in its operation.\nThe supply cassette 2 is preloaded with the wafer 14 provided with an ultraviolet sensitive tape 16 (herein after referred to as a UV tape) on the front surface as shown in FIG. 2 (tape laminated). The transfer machine 6 extracts the wafer 14 provided with the UV tape 16 from the supply cassette 2 and transfers it to the wafer grinder 8.\nThe wafer grinder 8 is generally formed of a wafer grinding section 18 and a drying section 20. First, as shown in FIG. 3, the transfer machine 6 fixes the wafer 14 on a chuck table 28 provided in the wafer grinding section 18. The chuck table 28 is connected to a vacuum device so as to hold the wafer 14 on the chuck table 28 by vacuum. The wafer 14 is fixed on the chuck table 28 with the UV tape 16 placed at the bottom.\nSubsequently, a grinding process (back grinding process) is performed until the wafer 14 has a required thickness. A grinding tool 22 comes into contact with the back surface (a surface with no element and provided with the UV tape) of the wafer 14 and then rotates to carry out this grinding process. In order to remove scraps and to cool the wafer 14, cleaning liquid is supplied to the surface subjected to the grinding process.\nAfter being subjected to the grinding process at the wafer grinding section 18, the wafer 14 is transferred to the drying section 20 by the transfer machine 6. Then as shown in FIG. 4, the cleaning liquid is dried using a light source 24 such as a halogen lamp.\nAs described above, after the back grinding process and the drying process on the wafer 14 in the wafer grinder 8, the transfer machine 6 extracts the wafer 14 from the wafer grinder 8 and then transfers the wafer 14 to the UV irradiation device 10. As shown in FIG. 5, the UV tape 16 is irradiated with ultraviolet (UV) rays from an ultraviolet (UV) lamp 26.\nA characteristic feature of an adhesive agent provided on the UV tape 16 is that it becomes stiff when irradiated with the ultraviolet rays. Thus by irradiating the UV tape 16 with the ultraviolet rays, the adhesive agent will become less adhesive. This facilitates the removal of the UV tape from the wafer 14.\nAs shown above, after irradiating the wafer 14 with the ultraviolet in the UV irradiation device 10, the transfer machine 6 extracts the wafer 14 from the UV irradiation device 10 and then transfers the wafer 14 to the remover 12. As shown in FIG. 6, the UV tape 16 is peeled off from the wafer 14 in the remover 12.\nIn the related art, the back grinding process is implemented on the wafer 14 using devices 2-12 and processes described above.\nHowever, in the related art, devices such as the transfer machine 6, the wafer grinder 8, the UV irradiation device 10 and the remover 12 were required separately in order to perform the back grinding process. Also, there was a need for the transfer machine 6 in order to mount and extract the water 14 between each of devices 6-12.\nAccordingly, it was necessary to provide a large number of devices 8-12 within a semiconductor manufacturing plant and to connect those devices 8-12 via a transfer machine 6. Therefore, there was a problem that a large installation space was needed within the plant.\nAlso, the transfer machine 6 mounts and extracts the wafer 14 between each device 2-12. Therefore, a considerable time is required between the extraction of the wafer 14 with the UV tape from the supply cassette 2 and the irradiation of the UV tape 16 at the UV irradiation device 10. This can cause a problem since the adhesiveness of the UV-sensitive adhesive agent provided on the UV tape may also change due to radiation from a fluorescent lamp.\nParticularly, when a considerably long time is required before ultraviolet irradiation of the UV tape 16 as in the related art, the property of the adhesive agent will be greatly changed by a fluorescent lamp. Accordingly, the adhesive agent will not be sufficiently stiff when the UV tape 16 is irradiated by the ultraviolet in the UV irradiation device 10. This causes a problem that the adhesive agent may remain on the wafer 14 when peeling off the UV tape 16 from the wafer 14 by the remover 12.\nAny residual adhesive agent on the wafer 14 will be inconvenient in the following manufacturing process for the semiconductor device. Therefore, in the related art, there was a need for a cleaning process after the removing process for removing any residual adhesive agent. This introduced an additional complexity to the process of manufacturing the semiconductor device."}
{"text": "There have been various attempts to make pieces of ice by forming an elongated section of ice and then somehow fracturing that ice section to form pieces which may or may not resemble ice \"cubes\", which are generally thought of as being formed individually, oftentimes in an individual cavity. Such apparatus is to be distinguished from continuous apparatus for making shaved ice or the like which operates on somewhat different principles and produces a soft ice.\nU.S. Pat. No. 3,984,996 discloses such an ice-maker which freezes a vertical tube of ice of square cross section and employs a central rotating shaft to drive the frozen tube upward. The apparatus is designed for continuous operation in a compartment of a home-freezer and employs a pair of synchronized blade-type cutters at the top to fracture the slowly upwardly moving tube of ice into cube-like sections. U.S. Pat. Nos. 2,595,588 and 3,287,927 also disclose ice-making machines which intermittently freeze a vertical column of ice and then eject that frozen column slowly upward where a breaking mechanism is provided for fracturing the emerging column into pieces at an upper location. None of these designs has proved to be entirely satisfactory, and it is not believed that any of them has been successful commercially. Accordingly, improved designs in ice-making equipment of this type have continued to be sought after."}
{"text": "The present invention relates to a single chain linear actuator and, in particular, to a single chain linear actuator having helically winding chains.\nOne measure of a linear actuator is the ratio of extended length to retracted length. For most linear actuators this ratio is 1.8 to 1. Although, there are linear actuators, such as scissor lifts and telescoping hydraulic jacks, in which the ratio of extended length to retracted length exceeds 3 to 1. This allows such linear actuators to provide a relatively large amount of movement while being positioned in a relatively small space. However, these linear actuators typically have a poor side load stiffness and/or require oil. For example, telescoping hydraulic jacks do not have particularly good side load stiffness and require oil. Accordingly, telescoping hydraulic jacks cannot be used in applications which require a good side load carrying capacity or where the use oil is not permissible.\nIt is known to provide linear actuators having extending and retracting chains, with inter-connected links, as disclosed in U.S. Pat. No. 5,271,182 to Greisner et al., and European Patent Application Number 1,484,465 to Soerensen. These prior art linear actuators have a good side load stiffness and do not require oil. However, they do not have a large extended length to retracted length ratio because a planar chain track, or chain path, is required within the actuator housing to guide the chain as it is being displaced. As a result, these types of linear actuators are typically used in applications where limited movement is required, or where there is no requirement for the linear actuator to be stored in a relatively small space."}
{"text": "The present invention is directed to an apparatus for packaging articles using shrink-wrap film and, particularly, to an improved heat tunnel that can be used for various film configurations.\nIt is known in the art to overwrap articles in a web of heat shrinkable film to form a multipack package by separating a tube of such film wrapped around spaced groups of articles along a weakened zone by shrinking the tube adjacent the zone and then by shrinking the tube section formed thereby around the articles to form a package. See U.S. Pat. No. 3,545,165.\nPrevious methods of packaging such as the above have involved feeding the groups of articles into a heat tunnel in series, with the film wrapped around the articles from the leading edge of the group to the trailing edge of the group. FIG. 1 shows how this is typically accomplished. Groups G of articles A are placed spaced apart on a conveyor C. A layer L of film F (usually from a roll of film) is wrapped around the groups G with the film layer L continuously covering adjacent groups G.\nThe groups G are then fed on the conveyor into a heat tunnel T. Heat and (typically) forced air is applied to the junction J between adjacent groups G, causing the film layer L to soften at the junction J and pinch off between the groups G, at the same time shrinking tightly against the groups G as shown. This results in complete packages P of articles A, with the film shrunk about them. The closed ends E of the packages P (known as “bulls eyes”) are at ends of the packages P in the direction of travel of the conveyor C (shown by the arrow).\nAn extension to the above apparatus is shown in FIG. 2. Here, parallel conveyors C1, C2, C3, etc. carry article groups G1, G2, G3, etc. into the heat tunnel T, where the above-described heat-shrinking occurs. The parallelism improves total throughput.\nThe apparatus shown in FIGS. 1 and 2 has a number of disadvantages. In gathering of multiple articles A into the groups G (known as “pack patterns”), the continuous tube of film F creates design challenges to support the groups G from the underside while the tube of film F is formed around the product. This is further complicated by product size changeover requirements. Theoretically, the conveyor C that transports the product pack pattern into the heat tunnel T would have to change widths for each change in product size to accommodate the tube of film F around the pack pattern.\nIn yet another variation (which the Assignee has used in the past), cut sleeves of film F are used, one sleeve per article group, instead of a continuous layer of film F around the groups G1, G2, G3, etc. However, the groups G are fed serially into the heat tunnel T with the articles A in each group G oriented in such a manner that the film F will be shrunk around each group G with the resulting closed ends E (“bulls eyes”) oriented transverse to the direction of travel of the conveyor C. To improve throughput, multiple parallel streams of articles A may be fed into the heat tunnel T.\nThe present application discloses an improved heat tunnel for use with both pre-perforated and non-perforated shrink wrap film.\nIn the packaging industry, aesthetics has become an increasingly important issue, both for the package that is produced and the machine that produces it. When the film is shrunk around the end of a package, it should leave a circular opening, the “bulls eye”, and should be free of wrinkles. This should be consistent from package to package and over a variety of product sizes.\nMany of the challenges in producing aesthetically pleasing “bulls eyes” stem from the way that current heat tunnels operate. Current heat tunnels often produce deformed bulls eyes due to uncontrolled airflow. That is, as the group of articles enclosed in shrink-wrap film enters the heat tunnel, the film is subjected to various disruptive air currents, causing the film to flutter as it is shrunk. This uncontrolled airflow results in the film wrinkling and shrinking non-uniformly, which in turn results in unaesthetically pleasing bulls eyes. Furthermore, current heat tunnels are not generally adjustable for various product sizes.\nThere is a need for a new heat tunnel capable of producing consistently good bulls eyes with controlled shrink and that is adjustable for a range of product sizes.\nThere is also a need for a new heat tunnel to reduce the heat transfer to the outer skin of the heat tunnel, increasing the operating efficiency and improving the working environment around the machine by lowering the temperature.\nThere is also a need for a more aesthetically appearing heat tunnel and one of reduced size.\nAll of the above needs are addressed by the present invention."}
{"text": "From a terminology point of view, a primary garment hanger is the hanger, be it metal, wood or plastic upon which a garment is placed for disposition upon a clothes bar or rod within the confines of a closet. A secondary garment hanger is a hanger which is to be disposed upon a clothes rod within a closet, and which is adapted to receive other hangers, which are in fact primary hangers as previously enunciated.\nMultigarment clothes hanger devices in the past have been multigarment primary hangers. Many men, for example, have metal or plastic devices having a plurality of spaced bars joined together to form a body which is connected in one fashion or another to a hook portion which in turn is disposed upon a clothes bar. While such devices are suitable for men's pants, they are not intended for suits or jackets, and in no way do they inhibit moth or other insect infestation.\nIt is known in the art to have plastic receptacles for paradichlorbenzene or moth balls. It is even known to provide hooks for mounting such receptacles upon clothes bars. While paradichlorobenzene does inhibit moth infestation, one can not wear clothes immediately after removing them from a closet which has been \"protected\" by paradichlorobenzene, due to the offensive odor of such chemical, and because some people are even allergic to the smell of moth balls. There is a need therefore for a clothes protector which can be brought into close proximity of a maximum of garments to afford insect protection, and which will not be malodorous.\nIt is an object therefore of this invention to provide a combination primary and secondary hanger device for use in clothing closets which is not malodorous.\nIt is another object of the invention to provide a secondary hanger which is adapted to inhibit clothes eating moth and other insect infestation in hanging garments for a large number of garments in the space normally occupied by about three garments on hangers.\nIt is a third object to provide a two piece hanger device that is capable of holding at least two and perhaps as many as twelve garments simultaneously, while saving closet space.\nIt is a still further object to provide a moth and other insect inhibiting hanger device that can be used as a primary and a secondary hanger either separately or simultaneously.\nOther objects of the invention will in part be obvious and will in part appear hereinafter.\nThe invention accordingly comprises the device which possesses the features properties and the relation of components which are exemplified in the following detailed disclosure and the scope of the application of which will be indicated in the appended claims.\nFor a fuller understanding of the nature and objects of the invention reference should be made to the following detailed description, taken in conjunction with the accompanying drawings."}
{"text": "This invention relates in general to the washing of truck or trailer body interiors, and in particular relates to apparatus and method for washing such interiors through an automatic washing cycle.\nVarious devices and systems have previously been employed for washing the interiors of truck or trailer bodies. Washing systems have been provided in which nozzles are mounted on small, wheeled vehicles which can roll into and out of the trailer. Such arrangements are relatively complicated and expensive to build and maintain, and furthermore are not fully automatic in operation such that the cycle time for washing a trailer is relatively long. Another washing system is of the type described in U.S. Pat. No. 3,689,318, issued Sept. 5, 1972, assigned to the same assignee as the present invention. The washing system which is disclosed therein employs a series of nozzles carried on the end of a boom which is fixed in position as a trailer being washed is backed and then driven forward so that the nozzles respectively enter and emerge from the trailer. In such a system the operator or driver controls the speed of the trailer as it is moved in relation to the nozzles so that the proper washing action is achieved."}
{"text": "This invention relates generally to the surveying of static structures, such as buildings, and more particularly surveying such structures for real estate tax assessment purposes."}
{"text": "A. Field of the Invention\nThe present invention relates to systems and methods for providing electronic communications to a customer. More particularly, the invention relates to systems and methods for providing an electronic account and other services to a customer by linking the customer's electronic address to a physical address where the customer receives physical mail.\nB. Description of the Related Art\nThe United States Postal Service (USPS) is an independent government agency that provides mail delivery and other services to the public. The USPS is widely recognized as a safe and reliable means for sending and receiving mail. With the steady growth of electronic communication and commerce, consumers and businesses need a secure way to communicate and conduct business electronically. Without trustworthy channels of communication, many potential participants in electronic commerce are unwilling to send sensitive information, e.g., credit card numbers, electronically, thus limiting the utility of electronic commerce to all individuals.\nElectronic mail, or e-mail, is a well-known means of communication for individuals and businesses with access to computers and Internet connections. When a user establishes an account with an e-mail service provider, e.g., America Online™ or Hotmail™, the user is assigned a unique e-mail address, e.g. joesmith@aol.com. Another individual can send a message to the user by entering the user's e-mail address along with the message and sending it via the Internet. E-mail can provide almost instant message delivery among individuals and businesses over vast distances for very little or no cost. E-mail also presents an opportunity for businesses to advertise to potential customers in a new way, e.g., by sending bulk advertisements via e-mail.\nDespite the advantages of e-mail, there are several drawbacks. Because e-mail is received and viewed electronically, e-mail does not reach those who are not “online.” In this way, e-mail contributes to the so-called “technology gap” between individuals with access to computers and computer technology and individuals who cannot afford or who do not understand computers and computer technology.\nAdditionally, the simplicity and low cost of e-mail make it an easy vehicle for unwanted messages, e.g., unsolicited advertisements or “spam.” Both individuals and businesses demand the capability to inhibit the receipt of unwanted e-mail.\nFurthermore, e-mail messages are also insecure, and can be intercepted en route by unknown third parties. Businesses and consumers who communicate electronically need to know that their messages are private, and that they can rely on the address to correctly identify the sender and/or recipient.\nTherefore, it is desirable to provide a system for communicating electronically that is available to everyone, that gives consumers control over the content of communications received, and that provides a secure and reliable way to conduct transactions electronically."}
{"text": "The present invention concerns a method for conducting a process in an automatically controlled system. A mathematical model of the process, with at least one variable model parameter, is implemented in a computing means. Before the process is started, the model precomputes at least one selected process parameter according to input values supplied to it for presetting the system. The input values and the process parameter are measured during the process and the precomputed value of the process parameter is adaptively improved after the process based on the measured input values. Furthermore, the present invention concerns a device for implementing a method for conducting a process.\nA method and a device for conducting a quasi-continuous process in an automatically controlled system are discussed in German Patent Application No. 40 40 360. Such processes typically include rolling trains where each pass of a rolled strip forms a process cycle (hereinafter, \"process run\"). Like all actual industrial processes, these process runs are time-variable. To conduct such processes, the system controlling the process must be preset before each run. That is, unlike traditional closed-loop control, control must precede the actual process because in industrial processes, controlled values can often be measured only indirectly and not at the point of affection of the process. Consequently, direct closed-loop control is not possible.\nThe system controlling the process is preset in a known manner by precomputing selected process parameters according to pre-established input values, or initially estimated input values, or both, based on a pool of relevant mathematical models of the process and by presetting the system using the precomputed parameters. Since mathematical models of the process to be conducted can only approximate the actual process, the model must be adaptively matched to the actual process events. For this, the process parameters and the input values are measured directly or indirectly by determining other measured values during each process run. After the process run is complete, the precomputation performed with the mathematical models is repeated within the framework of a postcomputation done based on the input values subsequently measured. The variable model parameters are adaptively modified based on the deviation between the process parameters thus computed and the measured process parameters to reduce the deviations obtained. The model parameters thus adapted are available at the beginning of the following process run for precomputing the process parameters.\nWith very complex relationships, partial aspects of the process to be modeled are usually described in partial models, and the interaction among the partial aspects of the process is described by a higher order model that links the partial models. While the partial models are still capable of describing the respective partial aspects of the process with sufficient accuracy, model assumptions for linking the partial models are difficult to make and such models may be subject to serious errors. In particular, adapting the higher-order model in addition to adapting the partial models is very difficult because the partial models provide no exact measured values, but only estimated values as input values for the higher-order model.\nThe goal of the present invention is to improve the precomputed value of the process parameter in the case of a very complex relationship between the process parameter and the input values."}
{"text": "In a wireless communication system designed to deliver data to a wireless mobile communication device (“mobile device”) such as a cellular telephone, a two-way pager, a wireless communication capable personal digital assistant (“PDA”), and other similar device, there are several main components in the wireless communication system. A host service, which provides services such as e-mail, calendar, and Internet web browsing, holds the data to be delivered to the mobile device. The host service is coupled to a router, which couples the host service and a wireless network that is designed to communicate with the mobile device. To make a timely delivery of the data, the host service forwards the data for the mobile device to the router when the data becomes available. The router then forwards the data to the wireless network, which transmits the data to the mobile device. If the mobile device fails to receive the data, the router or host service queues the data and re-forwards the data to the wireless network, which re-transmits the data to the mobile device. This process continues until the mobile device receives the data and acknowledges the reception or the process times out after a predetermined time period.\nA user of the mobile device is effectively logged onto a wireless service session when the host service and mobile device can establish a communication link by authenticating the logon credentials entered by the user. A session is deemed enabled active when the host service can authenticate the login and password entered. Thereafter, the device and host service can send and receive messages across the established wireless link.\nThe user typically logs onto a service session by entering a login and password at the beginning of the session (e.g, when user decides to start using the device and/or service provided by the host service). In some instances, this logon sequence (i.e., enter login/password, send to host service, authenticate with host service, establish connection) may be delayed, due to such scenarios as network. congestion or traffic. This may negatively impact the user's perception of the service or device to be slow and/or unreliable."}
{"text": "A fuel cell produces electrical energy by electrochemically oxidizing a fuel such as hydrogen and methanol in the cell to directly convert the chemical energy of the fuel into electrical energy. Fuel cells have recently drawn attention as a clean supply source for electrical energy.\nFuel cells are classified into a phosphoric acid type, a molten salt of a carbonic acid type, a solid oxide type and a solid polymer electrolyte type. Of these, the solid polymer electrolyte type fuel cell using a cation exchange mambrane as an electrolyte is called a proton exchange membrane type fuel cell. The proton exchange membrane type fuel cell is expected as a portable electric source such as an electric source for an electric car and a simple auxiliary electric source because it has high energy density even at a low operating temperature of 100.degree. C. or less.\nA proton exchange membrane type fuel cell comprises an ion exchange membrane and a pair of gas diffusion electrodes bonded to both sides of the ion exchange membrane. Each gas diffusion electrode has a catalyst at least on a side thereof facing the ion exchange membrane. The cell is operated by feeding a fuel such as hydrogen to one gas diffusion electrode and feeding an oxidizing agent such as oxygen and air to the other gas diffusion electrode respectively, and connecting an external load circuit to both gas diffusion electrodes.\nThat is, a proton (a hydrogen ion) and an electron are generated due to the oxidization of fuel at one gas diffusion electrode. The proton is transferred to the other gas diffusion electrode through the membrane by conduction and there, water is produced by the reaction of the proton with oxygen contained in the oxidizing agent. At this time, the electron generated at one gas diffusion electrode is transferred to the other through the external load circuit to obtain electrical energy.\nAs mentioned above, in a proton exchange membrane type fuel cell, an ion exchange membrane operates as an electrolyte to conduct the proton. Further, the ion exchange membrane substantially forms one body structure with the gas diffusion electrodes due to a bonding of the electrodes to both sides of the membrane. Therefore, the ion exchange membrane also plays a part of a diaphragm by not allowing fuel to mix directly with the oxidizing agent.\nThe ion exchange membrane used for the proton exchange membrane type fuel cell requires low electrical resistance, quick movement of water through the ion exchange membrane, high water retention characteristics to maintain low electrical resistance and permeability to gases, which allows oxygen gas and hydrogen gas to be fed to the electrodes at a high enough speed. In addition, the ion exchange membrane requires an appropriate permeability to gases, excellent chemical stability during prolonged use and strong physical strength in view of its role as a diaphragm.\nAs a conventional ion exchange membrane used for the proton exchange membrane type fuel cell, for example, NAFION (registered trademark) manufactured by E.I. du Pont de Nemours and Co. having a fluororesin as a main chain of a polymer and a sulfonic acid group as an ion exchange group is used.\nHowever, the conventional ion exchange membrane used for the proton exchange membrane type fuel cell can not respond to the increased request these days for a proton exchange membrane type fuel cell having high performance. The conventional ion exchange membrane is excellent in chemical permanence properties and stability. However, it has high electrical resistance. Further, it easily becomes dry due to low water retention characteristics so that proton conductivity is reduced or the reaction of fuel gas or oxidizing agent gas is inhibited at the electrode having a catalyst.\nInternational Unexamined Patent Publication No. Wo86/06879 discloses a diaphragm used for a fuel cell having an equivalent weight of less than 1000 g/eq and strong physical strength. This ion exchange membrane has relatively high strength even at a high temperature of 110.degree. C. or more. However, when it is used for fuel cell at a temperature of 100.degree. C. or less, its performance is not sufficient.\nEuropean Patent Unexamined Publication No. 0498076 discloses an ion exchange membrane having an equivalent weight of 700 to 1000 g/eq and a water content of 35 to 100% by weight. When it is used for a fuel cell, which is operated at a low pressure of about 1 atm., or uses air as an oxygen resource, its performance is not sufficient.\nNeither of the above-mentioned ion exchange membranes has necessary performance as a diaphragm and an electrolyte when used for a fuel cell.\nThe present invention has been completed to overcome the above-described problems of the prior art. That is, the present invention provides an ion exchange membrane used for a proton exchange membrane type fuel cell having excellent performance as a diaphragm and an electrolyte by specifying the molecular structure of the ion exchange membrane and limiting its electrical conductivity, permeability to gases and water content to appropriate ranges. The proton exchange membrane type fuel cell comprising the ion exchange membrane of the present invention can maintain high output performance for a long time."}
{"text": "Silicon (Si) nanowires, due to their potential of surpassing conventional lithography limitations and tuning their chemical and physical properties at the nanometer scale during growth, are attractive building blocks for high-performing functional electronic devices. While extensive research has been devoted to developing methods to assemble and integrate Si nanowires into planar functional devices, limited efforts have been devoted to the device integration of vertical Si wire arrays (both nanowires and microwires) even though such wire arrays are desirable for applications ranging from sensors, solar cells, thermoelectric devices, Li-ion batteries to vertical field-effect transistors.\nSi wire arrays with lengths of a few to tens of micrometers are usually formed by the deposition of vapor onto a Si wafer that is about 500 microns thick, rigid, and opaque, or by the chemical etching of such a wafer. Separation of the resulting Si wire arrays from the parent Si wafer allows the electrical, thermal, optical and mechanical properties of the Si wire arrays to be isolated and harnessed in devices without being overshadowed by the properties of the thick parent Si wafer, and allows the transfer of the Si wire arrays to other flexible, lightweight, low cost or transparent substrates for enhanced device functionality. A key condition for the separation and transfer of Si wire arrays is the preservation of their original properties and orientation. A range of techniques have been proposed for the separation of vertical Si wire arrays from their parent Si wafers, which rely on the mechanical breakage of the Si wires by the application of peeling forces, direct shear forces, or the creation of a horizontal porous crack within the Si wires. However, these techniques rely upon encapsulation of the Si wires in a polymer host so that the Si wires are held together during the mechanical breaking process, and the presence of the polymer host makes it difficult to use the separated Si wires for applications that benefit from an exposed Si surface, such as for Ohmic contacts to metal electrodes, sensors or catalysts. In addition, some of these techniques are restricted to Si wires of certain geometries or those made by particular methods.\nIt is against this background that a need arose to develop the electro-assisted technique described herein."}
{"text": "Communications systems often employ adaptive equalization to compensate for the distortion effects of changing channel conditions and disturbances on the signal transmission channel. The equalization process, for example, may estimate the transfer function of the transmission channel and apply the inverse of the transfer function to the received signal so as to reduce or eliminate distortion effects.\nChannel equalization typically employs filters that remove amplitude and phase distortions resulting from a frequency dependent time variant response of the transmission channel, for example, to thereby provide improved symbol detection capability. The channel equalization may remove baseband intersymbol interference (ISI) caused by transmission channel disturbances, including the low pass filtering effect of the transmission channel. ISI may cause the value of a given symbol to be distorted by the values of preceding and following symbols, and essentially represents symbol “ghosts” because ISI may include advanced and delayed symbols with respect to a reference symbol location in a given decision region.\nAn adaptive equalizer may be viewed as a digital filter with an adaptive response to compensate for channel distortions. Several well-known algorithms are available for adapting the filter coefficients and thereby the filter response to converge the equalizer.\nSignificant effort has been spent enhancing adaptation algorithms for use in data transmission, whether over communication systems or from storage mediums. Such adaptation algorithms typically are employed to compensate for distortions introduced into the signals by the transmission medium through which the signals have traveled. Such transmission mediums might comprise, for example, optical networks, wireless networks, standard public switch telephone networks, or even storage media where the signals have been stored and transmitted through interfaces to a user.\nIn all of these situations, the media through which the signal is transmitted or in which the media is stored may, in fact, affect the contents of the signal. Accordingly, equalization systems are generally employed to place the signal as nearly as possible in its original form. However, because the exact nature of this distortion as to which the signal is subjected to may not be known at the receiving end, the structure and methods employed to equalize the received signals necessarily involve certain assumptions.\nTypically, the techniques for enhancing the adaptation algorithms focus on the value of the signal at the time the signal is sampled in discrete time implementations. However, very limited effort has been spent with respect to continuous time adaptation for continuous time delay line equalizers.\nOne technique which has been used to adapt the feedforward equalizer to compensate for signal distortion due to the transmission or storage media is known as discrete time least mean square based adaptation. A typical continuous time least mean square (“LMS”) based adaptation is applied to the feedforward equalizer followed by decision feedback. This technique can be expressed as:\n                                          c            i                    ⁡                      (            t            )                          =                              ∫            0            t                    ⁢                      μ            ·                          e              ⁡                              (                τ                )                                      ·                                          s                i                            ⁡                              (                τ                )                                      ·                          ⅆ              τ                                                          (        1        )            where:                μ is an adaptation parameter;        ci denotes the ith tap either feed forward or feedback;        si(τ) denotes the sampled input signal appropriately time aligned as applied to the ith tap; and        e(τ) denotes an error signal, such as the difference between the output of a timing control circuit and a feedback signal.One problem with this form of the adaptation is “coefficient drift” associated with decision-directed fractionally spaced equalizers. For decision-directed equalizers, “coefficient drift” can be a serious problem for a number of reasons. For example, the coefficient vector c with coefficients ci satisfying the above equation can be c or −c. Furthermore, with enough taps, the feedforward equalizer coefficient vector can be shifted either left or right (i.e., the effective group delay change of the feedforward equalizer) with the same performance.        \nThe generic coefficient-update equation in continuous time LMSE equalizers can be given as follows:\n                                                        c              i                        ⁡                          (              t              )                                +                      T            ⁢                          ⅆ                              ⅆ                t                                      ⁢                                          c                i                            ⁡                              (                t                )                                                    =                              β            ⁢                                                  ⁢                                          c                i                            ⁡                              (                t                )                                              +                      μ            ·                          e              ⁡                              (                t                )                                      ·                                          s                i                            ⁡                              (                t                )                                                                        (        2        )            where T is the time-constant parameter, μ is the adaptation gain parameter and β is the leakage parameter. Note that when β=1 (representing no leakage), the equation simplifies to the previous integral equation. The coefficient-drift problem can be minimized with tap-leakage.\nIn general, a significant challenge within a communication system (e.g., a telecommunication network or a data retrieval system) has been the application of channel parameter measurement and performance monitoring techniques. These techniques are desired to assist with maintenance of the communication system and/or quality of service and possibly provide fault detection and isolation.\nFor example, the techniques may be utilized to provide relevant information or take remedial action during set-up or provide system optimization to enhance throughput, reliability, and/or monitor the health within the communication system on a real-time basis. It may also be important to have a measure of the quality of the signal to prevent any degradation in the quality of service through any path of the communication system.\nThus, adaptively changing the coefficients allows the system to compensate for the various changes in the signal, such as discussed above. Adaptive coefficient generators may be used to provide appropriate coefficients to a feedforward or feedback filter or equalizer. The adaptive coefficient generators changes the coefficients based on an input signal, such as an information or data signal s(t), and an error signal, which may be an iterative feedback error signal e(t). However, for LMS adaptation of coefficients, input signals s(t) and e(t) are ideally uncorrelated, i.e., the average value of the product is zero at the converged state, with optimal coefficients s(t) and e(t) uncorrelated. This would require infinite effective loop gain and hence translate to infinite gain requirement for the coefficient adaptation block. The high gain requirement poses a very severe constraint, e.g., resulting from mismatches in differential circuit topologies. Nominally identical devices suffer from random mismatches, and the problem worsens with the scaling of CMOS technology.\nOne possible way for correcting offsets is to use digital-to-analog converters (DACs). But, high device offsets, along with high gain, increases the range and resolution requirement of the DAC. Note that offsets need to be corrected at the power-up of the IC only, as in continuous-time mode, there is no way to distinguish between the adapted coefficient and the offset.\nAccordingly, it would be desirable to have systems and methods for adaptively generating coefficients for continuous time equalizers that overcome the disadvantages of the prior art as discussed, such as generating coefficients with high adaptation loop gain and requisite time-constant, mitigating offsets, and maintaining closed-loop stability."}
{"text": "1. Field of the Invention\nThis invention relates to a device for detecting the position of an image of an object, and more particularly to an image plane deviation amount detecting device of an automatic focusing apparatus in a camera.\n2. Description of the Prior Art\nA focus detecting device disclosed in U.S. Pat. No. 4,185,191 or U.S. Pat. No. 4,264,810 for making a first image and a second image by an imaging lens from light beams passed through two areas of the exit pupil of a picture-taking lens and detecting front focus, rear focus or in-focus from the amount of displacement of the first and second images relative to a pair of image position detecting photoelectric converters disposed in or near the focal plane of the imaging lens (hereinafter referred to as the amount of lateral deviation between the object images) is known in the prior art.\nIn a so-called automatic focusing apparatus for automatically rendering the picture-taking optical system into in-focus condition by the information of focus condition obtained from the focus detecting device of this type, the control of movement of the picture-taking lens as hereinafter described has heretofore been effected. By the information of front focus or rear focus from the focus detecting device, movement of the picture-taking lens toward the in-focus point is started. The focus detecting device detects the focus condition also during the movement of the picture-taking lens and stops the picture-taking lens when it detects the in-focus condition or the vicinity thereof.\nWhat has been described above can be regarded as the so-called feedback loop control using an optical system. However, where there is a response delay in the focus detecting device due to the charge accumulating time or the signal processing time when a light-receiving element such as CCD is used, if the picture-taking lens is moved at the usual speed of lens movement, there is the possibility that a situation arises in which the focus detecting operation cannot overtake the movement of the picture-taking lens. In such a case, there arises the so-called hunting in which the picture-taking lens goes past the in-focus point and is again moved back in the reverse direction and therefore, the lens movement must be done at a low speed and this may lead to the problem that the automatic focusing operation cannot be accomplished quickly."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of measuring pyrophosphate.\n2. Description of the Related Art\nAs conventional techniques of measuring pyrophosphate enzymatically, there are known, for example, (i) a method involving using pyruvate orthophosphate dikinase and luciferase (Japanese Patent Application Laid-Open No. 2007-097471), (ii) a method involving using hypoxanthine phosphoribosyltransferase and xanthine dehydrogenase/oxidase (Japanese Patent Application Laid-Open No. 2003-174900), and (iii) a method involving using kinase, an enzyme for producing ATP from pyrophosphate, and dehydrogenase which requires NAD or NADP as a coenzyme (Japanese Patent Application Laid-Open No. 2006-187251)."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for controlling a radiation image pickup apparatus suitable for medical diagnosis and industrial non-destructive inspection, the radiation image pickup apparatus, and a radiation image pickup system and, in particular, to a method for controlling a radiation image pickup apparatus, the radiation image pickup apparatus, and a radiation image pickup system capable of detecting the presence of irradiation, such as start and end of irradiation.\n2. Description of the Related Art\nA radiation image pickup apparatus including a flat panel detector (hereinafter simply referred to as an “FPD”) performs an image capturing operation in synchronization with irradiation performed by a radiation generating apparatus. To synchronize two operations, the following techniques are employed. That is, as described in U.S. patent application Ser. No. 2007/0125952, a radiation image pickup apparatus includes a plurality of pixels. Each of the pixels includes a conversion element that converts radiation or light into electric carriers and a switch element that transfers an electric signal based on the electric carriers so as to be capable of providing a desired voltage to one of two electrodes of the conversion element. The pixel further includes a detection unit for detecting the presence of irradiation. Before the detection unit detects irradiation from the radiation generating apparatus, a conducting voltage is sequentially provided from a drive circuit to the switch elements on a row-by-row basis in order to switch the switch elements to a conducting state. In this manner, the voltage of one of the electrodes of the conversion element is reset. Thereafter, if the detection unit detects start of irradiation from the radiation generating apparatus, supply of the conducting voltage from the drive circuit is stopped. Thus, a non-conducting voltage is supplied from the drive circuit to all of the switch elements in order to switch the switch elements to a non-conducting state. Accordingly, the electric carriers generated in the conversion element is accumulated in each of the pixels. In contrast, if the detection unit detects end of irradiation, a conducting voltage is sequentially supplied from the drive circuit to the switch elements on a row-by-row basis. Thus, an electric signal in accordance with the accumulated electric carriers is transferred from the pixel.\nJapanese Patent Laid-Open No. 2011-185622 describes that a noise component that is not caused by irradiation is coupled with a signal output from a detection unit for detecting irradiation due to a reset operation, and the noise component has a particular variation. Therefore, according to Japanese Patent Laid-Open No. 2011-185622, the profile of a signal used for detecting irradiation is measured in advance, and differential processing is performed on a measured signal for detecting irradiation using the profile. Thereafter, the signal subjected to the differential processing is compared with a predetermined threshold value to detect start of irradiation.\nAs described above, in order to accurately detect the presence/absence of irradiation, the threshold value needs to be optimally set. However, U.S. patent application Ser. No. 2007/0125952 and Japanese Patent Laid-Open No. 2011-185622 do not mention the setting of an optimal threshold value and, thus, further study is needed.\nAccordingly, the present invention provides a radiation image pickup apparatus capable of setting an optimal threshold value used for instantaneously and highly accurately detecting the presence/absence of irradiation."}
{"text": "This invention relates to water crafts and particularly to sailboats having keel structures for improving the lateral stability of the sailboats when pointing the sailboats. The lateral stability provided by the keel structure enables a sailboat to better hold its course.\nThe keel structure includes a keel member and a wing member attached thereto. In one embodiment the keel member is slidably attached to the sailboat, and the wing member is fixed to the keel member. In a second embodiment the keel member is pivotally attached to the sailboat, and the wing member is pivotally attached to the keel member.\nKeels are commonly provided on the bottom of sailboats to add stability to the sailboats by preventing the sailboats from heeling. Larger sailboats have been provided with fixed keels, while smaller sailboats have been provided with vertically moveable keels or dagger boards. While these fixed keels and dagger boards do increase the stability of the sailboats to which they are attached, the stability is not always sufficient to allow these sailboats to hold their course when pointing, and especially when pointing high, i.e. sailing close to the wind. Additionally, because of the fixed keels of the larger sailboats, these sailboats cannot be sailed in shallow water. Also, in deeper water the size of the fixed keels limits the speed of the sailboats. Whether sailing for recreational purposes or competitive purposes, the most enjoyment can be obtained when the sailboat is stable and versatile, such that it can be controlled as desired by its captain in changing water depths and wind conditions.\nIt is a primary object of the invention to improve the lateral stability of a sailboat when pointing the sailboat.\nIt is another object of the invention to enable a sailboat to easily be sailed in shallow water.\nIt is a further object of the invention to enable a sailboat to be sailed in deeper water at increased speeds.\nThe water craft of the present invention is preferably a sailboat provided with a wing keel. In a first embodiment of the invention the wing keel comprises a keel member that is pivotally attached to the bottom of the sailboat, and a wing member that is pivotally attached to the bottom of the keel member. As the sailboat is being sailed the force of the water will pivot the keel member towards the bottom of the sailboat, and also pivot the wing member such that it is generally parallel to the water surface. When the wing member pivots such that it is generally parallel to the water surface, a greater generally horizontal surface area is presented to the water, whereby the sailboat is provided with enhanced stability in all water depths and wind conditions while performing sailing maneuvers, such as pointing. Also, because the keel member pivots towards the bottom of the sailboat the draft is reduced, whereby the sailboat can be sailed in shallow water. This reduced draft also enables the sailboat to be sailed at greater speeds in deeper water.\nIn a second embodiment of the invention the sailboat includes a vertically adjustable keel member having a wing member fixed to the bottom thereof. The wing member is fixed to the bottom of the keel member such that the wing member extends generally parallel to the water surface during sailing of the sailboat, whereby stability of the sailboat is improved as with the first embodiment. And, because the keel member is vertically adjustable the draft can be reduced, whereby the sailboat can be sailed in shallow water and can also be sailed in deeper water at a greater speed, like in the first embodiment"}
{"text": "Modern microprocessors may support branch predictors in their architectures. In a pipelined architecture, the branch predictors may permit the fetch and execution of instructions subsequent to a branch point before the actual resolution of the conditional branch by execution of the branch instruction. This permits enhanced throughput when the branch predictor issues a correct prediction. There are many methods and hardware implementations for these branch predictors, but none of them are proof from errors. The enhanced throughput when the branch predictor issues a correct prediction must be weighed against the branch misprediction penalty that accrues when the branch predictor issues an incorrect prediction. The branch misprediction penalty may include such costs as stalling the pipeline during the branch misprediction recovery time and the execution of instructions along the mispredicted path that are potentially wasted.\nOften in programs the two paths (correct and incorrect) of a branch reconnect at a future point, which may be called a convergence point. It would appear that, if the branch was executed along a mispredicted path to a convergence point and beyond, the recovery could consist of re-executing only a portion of the instructions along correct path that were executed on the mispredicted path, and re-use the results of the rest of the instructions on the correct path that were executed on the mispredicted path subsequent to the convergence point. However, this could not generally be accomplished as the renamed physical registers could be contaminated with data from execution of the mispredicted path instructions between the mispredicted path and the convergence point."}
{"text": "Composite materials combine two or more distinct materials with complementary qualities, such as for instance lightness and strength. Various composite materials are known to the skilled person. For instance, honeycomb sandwiches, combining a honeycomb core and two facing panels, in metal, polymer and/or other materials, have long been used in a number of different applications, and in particular for structural elements in the aerospace and shipbuilding fields. Other composite materials combine a solid matrix of a first material with reinforcing elements, usually fibers, of a second material embedded in the matrix. Such composite materials include ceramic matrix composites (CMC), metal matrix composites (PMC) and polymer matrix composites (PMC). Advances in various fields, such as nanotechnology, have expanded the use of these materials to many technical fields, such as power generation, construction, medical implants and prostheses, transportation, etc. This has led to further competition to increase the performances and reduce the drawbacks of these materials.\nAmong composite materials, polymer matrix composites (PMC) and in particular fiber-reinforced polymers (FRP), such as, among others, carbon-, glass- and/or aramid-fiber reinforced polymers are particularly widespread. Fiber-reinforced polymers offer an advantageous combination of the properties, in particular the mechanical properties, of a polymer matrix and reinforcing fibers embedded in said polymer matrix. Both thermosetting and thermoplastic polymers are commonly used as matrices in such fiber-reinforced polymers. To produce a fiber-reinforced thermosetting polymer article, the fibers are first impregnated with a resin, i.e. a prepolymer in a soft solid or viscous state, shaped into a given form, usually by molding, and the resin is then irreversibly hardened by curing. During curing, the prepolymer molecules crosslink with each other to form a three-dimensional network. To initiate or at least accelerate this crosslinking reaction, the resin is usually energized using thermal heat transfer mechanisms and/or electromagnetic excitation. On the other hand, fiber-reinforced thermoplastic polymer composites can be produced by heating a thermoplastic so that it melts and impregnates reinforcing fibers. The production of fiber-reinforced thermoplastic polymer articles normally involves a heating stage in which the material is heated in order to soften the thermoplastic and enable processes such as forming or handling. Some preforms for fiber-reinforced thermoplastic polymer articles include a so-called comingled fabric in which the reinforcing fibers are mixed with thermoplastics. In this case the impregnation step takes place during forming.\nMicrowave heating technology is the most promising candidate for curing, drying, thermal treatment, inspection, post-consolidation, repair and a number of other processes for composite materials. A method for heating a fiber-reinforced polymer article using microwaves was disclosed in Japanese patent publication JP H5-79208 B2. According to this first prior art method, the fiber-reinforced polymer article is held in a mold made of a similar material with substantially the same dielectric properties. The mold containing the fiber-reinforced polymer is irradiated with microwaves, whose energy is converted into heat by both the mold and the fiber-reinforced polymer inside it. However, in this method, since the mold absorbs part of the microwave radiation, the dielectric heating of the fiber-reinforced polymer article may not be sufficiently homogeneous. In particular, in a thick-walled hollow article such as a pressure tank, the inner layers of the article could be insufficiently heated as a result.\nAnother method for heating up a fiber-reinforced polymer article using microwaves was disclosed in Japanese patent application Laid-Open JP H11-300766 A. According to this second prior art method, the fiber-reinforced polymer article is held in a mold made of a material that is substantially transparent to microwaves. In this method, the dielectric heating by the microwave radiation is substantially limited to the fiber-reinforced polymer, rather than the mold. However, this method also has the potential drawback of insufficiently homogeneous heating, in particular in thick-walled hollow articles.\nTo facilitate the penetration of the microwave radiation into the thick-walled article, it may be considered to dispense with the mold. However, when the fiber-reinforced polymer article has a high fiber-to-matrix ratio, a number of the embedded fibers, and in particular loose fiber ends, may be exposed. If the fibers are electrically conductive, as for instance carbon fibers are, the currents induced in the fibers by the microwave radiation during microwave curing will generate electric sparks and arcs, which can lead to local charring or even fire."}
{"text": "Cable television systems have made possible the transmission of many channels of television programs to the homes of subscribers. Instead of being limited to the number of VHF and UHF channels that can be transmitted and received in a given area, the number of channels in cable systems is limited only by the transmission characteristics of the cable itself and the ability to compress the information in television programs into narrower-bandwidth channels. The advent of fiber-optic transmission systems for use in cable television networks has also vastly increased the number of channels available.\nSuch increase in channels has also given rise to proposals for interactive television systems wherein a subscriber can transmit information or requests back into the system, which information or requests may subsequently affect programs or information directed to such subscriber. There are a wide variety of applications for interactive television systems, such as video games, video catalog shopping, teaching systems, movies on demand and audio programs. Each application can be tailored for an individual subscriber, for example, a subscriber may be able to select the language of the soundtrack in a movie. However, such systems typically require the ability (i) to control specific programs or information sent to each subscriber and (ii) to receive input messages or requests from the subscriber.\nIn many cable television systems, a special converter is used at the subscriber location to allow the subscriber to select among the various available channels, and possibly to \"unscramble\" premium channels for which extra tees are payable. Such converters usually provide an output on one of the channels that a standard television receiver can receive, such as Channel 3. More recently, television receivers and video recorders have been made available that are \"cable ready,\" that is, with the ability to receive and select among all the channels transmitted over the cable system, but without the ability to unscramble premium channels. A converter is still needed for the premium channels.\nIn order to implement many of the features of interactive television systems, it is necessary for messages containing information or requests from subscribers to be transmitted to a central control location. It is desirable for these \"uplink\" messages to be transmitted from subscribers over the same facility used to distribute programs. Conventional cable systems use frequencies above 50 MHz, allocated in 6 MHz distribution channels, for programs and frequencies in the 5-30 MHz band for such uplink messages. However, because many subscribers will be sending messages over the same uplink frequencies, some means is needed to coordinate such messages to prevent interference and to identify the source of each message.\nWhen a subscriber is first connected to a cable television system with converters at subscriber locations, it is typically necessary for a service technician to install and initialize a converter for such subscriber in the field. Also, a visit by a technician may be needed to re-initialize a converter after a service interruption or a power failure. Such visits can be costly. Accordingly, it is desirable to provide a converter that can be initialized remotely over the cable system so that such a converter can be furnished to a subscriber for self-installation and for remote initialization or re-initialization."}
{"text": "1. Field of the Invention\nThe present invention is directed generally to the field signal processing and, more particularly, to the field of video signal processing.\n2. Description of the Background\nIt is known in the art to use color as a basis for merging pixel by pixel, two video signals. Chroma-keying, as the process is known, uses color information to control a switching process. A person, for example a weatherperson, stands in front of a blue screen while a first video signal is made. A second video signal is provided which is to be combined with just the weatherperson from the first video signal. The first video signal is analyzed pixel by pixel for the key, i.e., the blue of the background screen. For each blue pixel, the corresponding pixel from the second video signal is selected for the final image. Each non-blue pixel from the first video signal is selected for the final image. Through that simple process, the weatherperson is extracted from the background screen and superimposed on top of the images in the second video signal. The chroma-key process is simple enough to be performed in real time. However, the extracted image is always considered to be in front of the images from the second video signal. That may or may not result in a realistic looking final product. Furthermore, because the extracted image is always in front of the images from the second video signal, there is no physical or geometrical interaction allowed between the extracted image and the images of the second video signal.\nOther methods for combining video signals are disclosed in U.S. Pat. No. 5,400,080 entitled Apparatus And Method For Combining Video Signals Representing Images Having Different Depths, U.S. Pat. No. 5,353,068 entitled Video Signal Combining Apparatus And Method, and U.S. Pat. No. 5,280,337 entitled Depth-Based Video Combining. Each of those patents discloses a method which uses depth information for combining one or more video signals. However, the patents do not disclose how to obtain the needed depth information, which is a non-trivial problem if pixel by pixel depth information is to be provided in real time for standard video signals having a frame rate of thirty frames per second. Moreover, there is no recognition in the aforementioned patents of how images from the combined video signals might interact, e.g. one image might cast a shadow on an image which is behind it in the combined video. Thus, the need exists for a method of using depth as the key for merging two video signals in real time in a way which enables the images from the combined video to interact with one another."}
{"text": "1. Field of the Invention\nThe present invention relates to a way of structuring input/output plugs and a digital interface for equipment in a communications system where a number of equipment are capable of multiplexing and transferring packets of control signals and information signals via a bus.\n2. Description of the Related Art\nIn conventional communication systems where AV (Audio/Video) equipment such as, for example, Video Tape Recorders (hereinafter referred to as VTRs), Televisions (hereinafter referred to as TVs) etc. are connected together via analog AV signal lines and digital control signal lines, a D2B (Domestic Digital Bus) has been used. Examples of this kind of system are described in the European patents that belong to the same applicant relating to this kind of communications system.\nFirst, a description will be given of an example of this kind of communications system with reference to FIG. 17. This communications system comprises a TV, a VTR 1, a VTR 2, a Multi Disc Player (hereinafter referred to as an MDP), and an editor. The VTR 2 is also connected to a device, which is not compliant with a digital control signal line (hereinafter referred to as a non-bus compliant device).\nThe AV signal input/output plug for the non-bus compliant device is connected to input/output plugs of other devices via analog AV signal lines, wherein the plugs are marked only with designation numbers P1, P2 and P3, protruding directly out from units known as switchboxes (SW box). Commands such as connection control commands are sent and received by other devices via analog AV signal lines and digital control signal lines, which are connected separately from AV signal lines. Each of the equipment is also comprised of one or more functional units, which in the case of a VTR would be (i) a deck for recording and playing back and (ii) a tuner for selecting the signal to be received, and in the case of a TV, would be (i) a monitor and (ii) an amplifier. There is also an AVC (Audio/Video Controller) for controlling the operation of all of the items of equipment, although this is not shown in the diagrams. In the following description, these functional units may also be referred to as sub-devices.\nControl connection can be carried out in two ways in a communications system structured in this way. These two ways, which will be referred to as connection control method 1 and connection control method 2, will be described below in numerical order.\nIn connection control method 1, information pertaining to the connection structure for each device, i.e. which plug is connected to an opposite equipment and whether plugs are used for input or output, is pre-set by the user. In this way, when the equipment is receiving the connection control command, the sub-device which is to be the AV signal source is connected to an output plug being considered as the destination. Alternatively, a path may be set up within the device between the input plug of the designated number and the appropriate output plug, then a command is transmitted to the device connected to the output plug. The object is then achieved at the point in time when the command reaches the sub-device within the specific device indicated by the destination. At this time, commands are not transmitted via the control signal line when the device is a non-bus compliant device. The plug numbers for devices connected to this kind of device have to be directly specified (for example, in the case in FIG. 17, the input plug P1 for the VTR 2 is specified).\nThe following is a description with reference to FIG. 18 of the sequence in connection control method 1 when an MDP output is recorded on a VTR 2 as a result of instructions from an editor. FIG. 18(a) shows the communications sequence and FIG. 18(b) shows the commands.\nFirst, the MDP receives a command giving an instruction to connect the deck output to the deck of the VTR 2 from the command master (editor) (communication 1). At this time, none of the output plugs of the MDP are connected to the VTR 2. However, the output plug P3 is connected to the input plug P1 of the VTR 2, so a command is sent to the VTR 1 via a path from the input plug P1 to the switchbox (communication 2).\nWhen the VTR 1 receives this command, none of the output plugs are connected to the VTR 2 but the output plug P3 is connected to the input plug P1 of the TV. A command is sent via a path from the input plug P1 via the switchbox to the TV (communication 3).\nWhen the TV receives this command, it is understood that the output plug P3 is connected to the input plug P2 of the VTR 2 and a command is sent to the VTR 2 from the input plug P2 via the deck (communication 4).\nThe VTR 2 then receives this command and switches the switchbox connection so as to input from the input plug P2 to the deck. Once the VTR 2 completes the process 4, the TV is notified of the process completion. Upon receiving this notification, the TV gives notification to the VTR 1 that the process 3 is completed. When the VTR 1 then receives the notification, it gives notice to the MDP that the process 2 has been completed. Upon receiving this notification, the MDP gives notice to the command master that the process is completed. These packets are omitted from the diagrams.\nHere, an editor has been set up as the command master. However, rather than setting an editor up on the communications system, a structure may be adopted where the MDP and the VTR have command master functions.\nIn a second connection control method, an equipment being central (in the following, this is an AV center) manages all of the connection information regarding which kinds of equipment are connected together via which numbered ports in which directions. If the AV center receives a command requesting AV signal connection via the digital control signal line, it sends this command via the digital control signal line to the target device. The target device then receives this command and performs the input/output selection. At this time, it is possible to designate the sub-device by category via the initial connection request from the command master (by BS, CS etc. which are more in line with the objectives, rather than by a tuner etc.). However, at the time of a connection request, only designation of its own plug is available because of it is only aware of its own plug structure.\nThe following is a description with reference to FIG. 19 of the sequence in a connection control method 2 when a TV is taken-as an AV center and output from an MDP is recorded by a VTR 2. Here, FIG. 19(a) shows the communications sequence and FIG. 19(b) shows the commands.\nFirst, the TV receives a command from the command master instructing it to connect the MDP deck output to the deck of the VTR 2 (communication 1). The TV then receives this command and sends a command to the MDP to the effect that the output for the deck sub-device is to be outputted from the output plug P3 (communication 2).\nIf the MDP receives this command, the switchbox is switched so that the deck output is connected to the output plug P3 and notification of the completion of the switching process is then sent to the TV. If the TV then receives this information, a command is sent from the input plug P1 via the output plug P3 to the VTR 1 (communication 3).\nIf the VTR 1 receives this command, the switchbox is switched so that a connection is made from the input plug P1 to the output plug P3 and notification of completion of the switching process is sent to the TV. If the TV then receives this notification, a command is sent from the input plug P2 to the deck to be inputted to the VTR 2 (communication 4).\nProviding the VTR 2 receives this command, the switchbox is switched so as to input from the input plug P2 to the deck. The TV is then notified when the switching process is complete.\nIf the TV then receives notification of completion from the VTR 2, notification of completion of the connection is sent to the command master.\nHowever, with connection control method 1, the digital control signal line becomes congested as a result of commands being transmitted between the equipment, which neighbor each other. An infinite loop may be formed depending on the settings of the equipment in the connection line. It is also necessary to know the structure of the equipment, which command master intends to control and the structure of the connections among the equipment in the overall system, so as to designate the plug numbers.\nAlso, in connection control method 2, no matter how simple a connection is made between a plug and the equipment to which the plug is connected, this may not be achieved without a request being made to the AV center. Further, designation can be made using the sub-device category at the time of a connection request, but plug designation can only be achieved for the equipment's own plugs, for which the structure is known.\nIt is therefore apparent that connection control of a variety of AV signals by holding connection information for input/output-plugs as interfaces with external equipment is not ideal. It is therefore the object of the present invention to provide a communications system capable of resolving these kinds of problems."}
{"text": "1. Field of the Invention\nThis invention relates to acoustic absorption and damping materials, and more particularly, to acoustic absorption and damping materials that utilize a piezoelectric phenomenon to convert mechanical energy into electrical energy and to subsequently dissipate the converted energy as heat.\n2. Description of Related Art\nAbsorbing or damping unwanted acoustic or vibrational energy involves converting that energy into another form, usually heat. At the molecular level, the only distinction between heat energy and acoustic or vibrational energy is the randomness of the vector directions of molecular displacements. Acoustic and vibrational energy is highly correlated with large numbers of molecules displacing at the same time and in the same direction. Heat in a particular object may well have the same or more energy than propagating acoustic or vibrational energy, but the motion of the molecules is random with the mean molecular displacement at any given location being near zero.\nTwo primary techniques are available for randomizing the vector directions of the molecules in a matrix material propagating acoustic or vibrational energy. Cushman, et al. (U.S. Pat. No. 5,400,296) teach the use of two or more species of particles with differing characteristic impedances in a matrix material to promote random internal reflections at boundaries within the matrix material and the subsequent increase in probability that phase cancellation at adjacent or nearby locales can take place. Single particle species may also be used in this manner, but with less effect. Phase cancellation effectively randomizes the vector direction of molecular movement where it occurs. A second approach involves the careful choice of materials that exhibit a high degree of internal hysteresis. This internal hysteresis is thought to be caused by metastable molecular energy levels within the material. Propagating acoustic or vibrational energy may boost a particular molecule into a higher energy level, thus subtracting that energy from propagating energy, where the molecule remains for some time before randomly returning to its original energy level. For a discussion of this effect see Hartmann and Jarzynski, \"Ultrasonic hysteresis absorption in polymers,\" J. Appl. Phys., Vol. 43 , No. 11, November 1972, 4304-4312.\nInstead of randomizing molecular displacements to dissipate propagating acoustic or vibrational energy, some of this energy can be removed by converting the mechanical energy of sound or vibration into electrical energy utilizing the piezoelectric effect. A piezoelectric material such as polyvinylidene fluoride (PVDF) may be polarized and a coating of a conductive material such as aluminum applied to produce a piezoelectric transducer that will convert acoustic energy into electric energy, thus facilitating removal of converted energy from the system. This approach is reported in a recent issue of the Japan New Materials Report (May-June, 1995, p 9). In this report acoustic energy reductions of up to 90% are claimed in material specimens only 10 to 30 microns thick. However, the need to polarize the material and apply conductive electrodes to tap off the electrical energy produced limits the usefulness of this technique."}
{"text": "Tissue products are in almost constant use in daily life. Toilet tissue, facial tissue, tissue wipes, and paper towels are examples of tissue products used throughout home and industry. Tissue products can be made of a single lamina (or ply) or can be laminates (or multi-ply, or else laminated) formed from two or more plies. As used herein a “lamina”, or “ply”, refers to a single sheet of tissue paper, and the term “laminate”, or “multi-ply”, or “laminated”, is used to characterize a paper made by uniting several layers of single sheets or plies together or to another woven or nonwoven substrate to yield a unitized material. The phrase “tissue laminate”, as used herein, refers to a laminate including at least one ply of tissue paper unitized with another substrate or ply, and these are well known and recognized by those skilled in the art. “Tissue paper” or “tissue” as used herein, refers to woven or nonwoven substrates that are relatively low weight sheets for use as sanitary products such as facial and bathroom tissues or paper towels. These also may be fabricated in condenser, carbonized, wrapping and cleaning grades. General characteristics may include one or more of softness, adequate strength and absorbency, clean appearance, and freedom from coarse and abrasive particles. See, e.g., G. A. Smook (1990), Handbook of Pulp & Paper Terminology; A Guide to Industrial and Technological Usage, chap. 16; Angus Wilde Publications, Vancouver, BC. ISBN 0-9694628-0-8. In typical embodiments, a tissue may have a sheet basis weight of up to about 50 grams per square meter, more typically about 13 to about 46 grams per square meter.\nA single ply tissue product has several drawbacks. For example, a single lamina tissue product will be stiffer than a dual laminae tissue product having the same total basis weight. This increased stiffness results in a consumer perception that the single lamina tissue product is not as soft as a multi-ply tissue product and may, therefore, be less preferred than the multi-ply tissue product. A laminated tissue having the same total basis weight as a single lamina tissue product can be made to have greater caliper. This increased caliper results in the consumer perception that the product has thickness (bulkiness) and high quality. Also, to provide consumers with the convenience of using a predetermined length of material, tissue laminates may be perforated at pre-selected lengths.\nHowever, tissue laminates are subject to the phenomenon of skinning. Skinning occurs when the laminae separate from one another and no longer remain intact to form a unitary laminate. Skinning may occur, for example, when trying to reduce tissue product packaging and transportation costs. For example, U.S. Pat. No. 4,886,167 discloses packages of toilet tissue and paper towels compressed orthogonal to the cores to reduce the core volume. This compression of the tissue product may impart shear forces to the laminae and result in skinning. Skinning may also occur during a tissue converting operation when the top sheet of a roll breaks off as a result of friction in the area of contact between the surface in the production line and the roll of paper. It can also occur during handling by the consumer. In the case of perforated toilet tissue skinning can lead to mismatched plies along the perforation lines when a portion of the top sheet of a roll is torn or lost.\nSeveral attempts have been made in the art to join tissue laminae in a manner to reduce or minimize skinning. One approach involves the application of mechanical forces (friction and compression) to interlock the sheets and keep them from separating from each other. These forces are generated by piercing pinholes in the laminates, with the edges of the knurls engraved in a steel roll, to marry the sheets, or by pressing wavy lines in the tissues during embossing. Another approach uses an adhesive to bond laminae together. This approach has been described, for example, in U.S. Pat. Nos. 5,143,776; 4,885,202; 4,806,418; 4,770,920; 4,481,243; and 4,513,051.\nSome special, water dispersible or soluble hot-melt adhesives, like the ones described in U.S. Pat. Nos. 6,448,463B1; 6,087,550; 6,103,809; and 5,663,286 and in WO99/35189, all assigned to H. B. Fuller Company, can be used. Tissue laminates made with water-soluble adhesives readily break down upon wetting, as, for example, might occur when toilet tissue or facial tissue is disposed of in a sewer system. Consequently, the adhesive used to bond tissue laminae together is desirably aqueous-based (solution, dispersion, emulsion, latex or the like) and at least water-soluble (solution) or water dispersible to some degree.\nNotwithstanding the success of hot-melt adhesives, it may be desirable to use aqueous-based adhesives for economical reasons. However, using aqueous-based adhesives to manufacture tissue laminates is problematic due, at least in part, to the water transmissivity and delicate nature of the individual tissue plies when wetted. Water-based formulations tend to bleed through (or strike through) the ply(ies) to which the formulations are applied. Bleed through leads to unplanned bonding among laminae beyond the desired design of the laminate, for instance. During production, efforts are made to minimize or eliminate this undesired phenomenon. Such efforts include reducing the amount of adhesive used as well as pre-wetting of the plies. These kinds of efforts can reduce bond strength between plies.\nMoreover, when coated onto a delicate tissue substrate, the tissue is susceptible to damage when impacted by adhesive at conventional application pressures. The water in the adhesive composition also can distort or even destroy aesthetic embossing. Wet plies are also more susceptible to damage, complicating manufacture. Drying also takes a relatively long time, decreasing productivity, when using aqueous adhesives.\nAccordingly, the tissue industry could benefit from a better way to use aqueous adhesives to manufacture delicate tissue laminates that would provide the advantages offered by such compositions while reducing or even eliminating one or more of bleeding through, undue loss of bonding strength, distortion or destruction of embossed or other three-dimensional features, loss of strength during manufacture when plies are wetted, and relatively long drying times. It would also be advantageous to apply a process that would allow for waste reduction of the adhesive during manufacturing of the laminated tissues as well as the reduction of the adhesive consumption. These benefits, particularly waste reduction, are also applicable to non-tissue products where lamination and/or bonding are desirable."}
{"text": "The present invention relates to a low temperature chemical vapor deposition process for producing indium-containing semiconductor materials employing triisopropylindium as the source of indium.\nIndium containing semiconductor materials have been grown by organometallic vapor phase epitaxy (OMVPE) using conventional organoindium precursors such as trimethylindium, triethylindium, ethyldimethylindium and cyclopentadienylindium. The most common OMVPE indium precursor presently employed is trimethylindium. However, there are several problems associated with the use of trimethylindium. Since trimethylindium is a solid at or below room temperature, the effective vapor pressure of trimethylindium in a bubbler changes with time due to changes in the surface area of solid. This results in transport problems in a OMVPE system. Since trimethylindium decomposes slowly at temperatures lower than 400.degree. C. in an atmospheric reactor, semiconductor film growth below 400.degree. C. is typically inefficient because trimethylindium is too stable at these low film growth temperatures. Furthermore, the use of trimethylindium in low temperature film growth can result in significant unintentional carbon impurity incorporation--due to the formation of methyl radicals during pyrolysis--which can be deleterious to the performance of a semiconductor device.\nAdditionally, the incorporation of Bi into InAs, useful for far infrared device applications, requires growth temperatures as low as 275.degree. C. At this low growth temperature, trimethylindium decomposition is not complete so the growth rates of InAs and InAsBi are too low.\nWhile ethyldimethylindium is a liquid, its relative compositional stability as a pure molecular compound is questionable. Furthermore, ethyldimethylindium is too stable for low temperature film growth and its use can also result in unintentional carbon incorporation due to the formation of methyl radicals during pyrolysis. Although triethylindium is a liquid, it has been reported to be marginally stable and has been observed to decompose in the \"bottle\" during storage. Furthermore, triethylindium undergoes parasitic reactions with Group V hydrides such as AsH.sub.3 and PH.sub.3 to form nonvolatile adducts upstream of the substrate in the atmospheric pressure OMVPE reactors. Cyclopentadienylindium is a solid with a vapor pressure an order of magnitude less than trimethylindium and appears to be thermally more stable than trimethylindium which rules out its use for low temperature film growth of indium-containing semiconductor materials.\nTherefore, there is a need for an alternative organoindium precursor without methyl groups or substituents that result in unintentional carbon incorporation for the low temperature chemical vapor deposition of indium-containing semiconductor materials. The availability of alternative chemical vapor deposition organoindium precursors with lower pyrolysis temperatures and lower carbon incorporation characteristics would greatly enhance the development of indium containing semiconductor materials and devices.\nOne object of the present invention is the provision of an improved process for forming indium-containing semiconductor materials using chemical vapor deposition.\nAnother object is to provide such process employing an indium precursor with a lower pyrolysis temperature than organoindium precursors heretofore used.\nStill another object is the employment in the above process of an indium precursor that is more efficiently pyrolyzed at lower film growth temperatures than organoindium precursors heretofore used.\nA still further object is to provide a process of the above type employing an organoindium precursor that results in less unintentional carbon impurity incorporation at lower film growth temperatures than with the use of conventional organoindium precursors such as trimethylindium and ethyldimethylindium precursors.\nAnother object is the provision of an improved indium-containing semiconductor material obtained by the invention process.\nOther objects and advantages of the invention will appear hereinafter."}
{"text": "Turbomachines, such as stationary gas turbines or aircraft engines, include, inter alia, a plurality of blades rotatably mounted on a rotor so as to either compress the fluid flowing through the turbomachine or to be driven into rotation by the fluid.\nIn order to minimize flow losses between the rotating blades and a surrounding flow duct boundary, the gap between the blade tips at the radial ends of the blades and the flow duct boundary must be as small as possible in order that as little fluid as possible may flow through the gap between the flow duct boundary and the blades.\nFor this reason, in known turbomachines, so-called labyrinth seals are provided for sealing between the blade tips and the flow duct boundary. In such labyrinth seals, the blade tips move in a groove which is formed by the blade tips cutting into a sealing material on the flow duct boundary during operation of the turbomachine. In this connection, it is also known to provide the blade tips with so-called blade tip hardfacings. Blade tip hardfacings have particles of hard material embedded in a metal matrix, which are intended to cut the groove for the labyrinth seal into the opposite sealing material of the flow duct boundary and to protect the blade tip from wear.\nMoreover, blades of turbomachines also include protective coatings, such as erosion-protection coatings, on the airfoil to also protect the blade material in the airfoil region from wear caused, for example, by erosion.\nAccordingly, a blade of a turbomachine must be provided with different coatings adjacent to one another, such as, for example, a hardfacing on the tip and an erosion-protection coating on the airfoil. However, care must be taken that none of the coatings, especially the blade tip hardfacing, becomes covered by the other coating, namely the erosion-protection coating, because the covered coating; i.e., for example, the blade tip hardfacing, is otherwise unable to perform its function in the desired manner. In fact, in the event of a failure of the cutting function of the blade tip hardfacing, the blade can be massively damaged by excessive thermal and mechanical loading.\nAccordingly, the individual coatings must be suitably applied one after another while preventing unwanted mutual coverage of the coatings.\nDE 10 2010 049 398 A1 describes a wear- and oxidation resistant turbine blade including an oxidation-resistant metallic coating, in particular an MCrAlY coating, in which M is a metal, in particular nickel, cobalt, or a combination thereof, and which, in addition, may have a ceramic thermal barrier coating. In addition to this oxidation-resistant protective coating, a protective coating of abrasive material and binder material is applied to the blade tip by laser metal forming. Initially, the oxidation-resistant protective coating in the form of the MCrAlY coating is applied to the blade over the entire surface thereof. Then, the MCrAlY coating is mechanically removed in the region of the blade tip, and subsequently the wear-resistant protective coating is applied to the blade tip in the blade tip region by laser metal forming."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a semiconductor device structure and to a method of manufacturing the same, and more particularly, the present invention relates to a transistor having a hetero grain stack gate structure, and to a method of manufacturing the same.\n2. Description of Related Art\nSemiconductor devices such as field effect transistors (FETs) are becoming increasingly important in low voltage applications. As semiconductor devices are fabricated to have a higher degree of integration, a faster operating speed, and a lower power consumption, the magnitude or size of a complementary metal-on-semiconductor field effect transistor (CMOSFET) included in the device is rapidly reduced. As FET devices are scaled to smaller and smaller dimensions, manufacturers must refine transistor designs to maintain optimum device performance.\nTo control problems associated with a short-channel-effect (SCE) in a deep-submicron MOS transistor, a dual gate type CMOSFET design is widely used. The dual gate type CMOSFET has a surface channel in each of the NMOS and PMOS transistors, and a symmetric low threshold voltage (Vth). For example, with a currently-available conventional dual gate device, when the threshold voltage of the NMOS device is 0.5V, then the threshold voltage of the PMOS device is −0.5V. The dual gate type CMOSFET uses N+ and P+ polycrystalline (polysilicon) gate electrodes in the NMOS and PMOS transistors, respectively.\nA conventional CMOSFET structure, and relevant steps of a method of manufacturing such a device, will be described with reference to FIGS. 1A-C as follows.\nReferring to FIGS. 1A and B, a standard twin retrograded-well process is used to form an N-well 2n and a P-well 2p in the substrate 2 via a conventional photoresist masking and ion implantation process. In the N-well 2n and P-well 2p are formed lightly doped drain regions 16a and 74a, including halo doping profiles 16b and 74b respectively. On each of the N-well 2n and P-well 2p is formed a gate oxide layer 6p and a stacked gate electrode on top of the gate oxide layer 6p, including a thin polysilicon layer 8p and stacked polysilicon structure 73. Laterally adjacent to the gate electrode is a liquid phase deposition (LPD) oxide layer 18 and stacked polysilicon layers 72. Such a structure is then subjected to a heavy ion implantation using, for example, BF2+.\nAs shown in FIG. 1B, after doping for both the PMOS and NMOS transistors is complete, a thermal treatment is performed to condense then LPD oxide 18 and to activate the S/D implants. The thermal treatment diffuses the heavy implants in the polysilicon stack layers 72 into the substrate 2 to form the buried contacts and ultra-shallow junctions. Then a refractory metal layer 28, such as Ti, Co, W etc., is sputtered on the polysilicon stack layers 72 and the LPD oxide 18.\nTurning to FIG. 1C, next a standard two-step silicidation process is performed, consisting of a rapid thermal annealing (RTA) to react the metal layer 28 with the polysilicon stack layer 72 to form a thin silicide layer 28a on the polysilicon stack layers 72. Then the unreacted metal is removed from the oxide using a standard wet etching process. A second RTA process is performed to transform the silicide's phase to a less resistive phase.\nHowever, there are some problems with this conventional CMOSFET device.\nIn the PMOS transistor, when the thickness of the gate oxide layer 6p is less than 50 Å, boron implanted into the polysilicon gate electrode diffuses into the gate oxide layer 6p (Boron penetration). When this happens, the boron diffuses into the semiconductor substrate 2 and decreases the carrier mobility which causes a threshold voltage fluctuation. Due to the threshold voltage fluctuation, the threshold voltage Vth of the gate cannot be controlled, thereby degrading the characteristics of the PMOS transistor.\nAlso, even though the polysilicon gate electrode is over-implanted by ions, it is not a complete conductor. Therefore, during the operation of the MOS transistor, a depletion region may arise due to a depletion of the electric charges at an interface between the gate oxide layer and substrate. The depletion region of the polysilicon gate has a magnitude of several angstroms (Å) and acts as a connected gate oxide layer capacitor. If the thickness of gate oxide layer is small, the characteristics of the transistor are poor due to the polysilicon-gate-depletion region.\nTo address the boron penetration and the polysilicon-gate-depletion effect (PDE) problems, a method of replacing a conventional polycrystalline silicon (poly-Si) gate with a polycrystalline silicon germanium (poly-SiGe) gate has been proposed (see, e.g., IEDM Tech. Dig. 1990 pp. 253-256). The poly-SiGe gate structure can be manufactured by the conventional CMOS process and can control the boron penetration and the PDE problems discussed above.\nAlso, the work function of poly-SiGe is different from that of poly-Si, thereby increasing the threshold voltage of the device. Therefore, the amount of channel doping can be decreased, increasing the carrier mobility.\nHowever, when the metal silicide layer of Ti or Co is formed on the gate structure (as described above with respect to FIGS. 1A-C), it is difficult to form it on the poly-SiGe layer compared with forming it on the poly-Si layer. Also, due to the penetration of Ge into the silicide layer, the resistivity of the poly-Si layer is rapidly increased, which is undesirable.\nTherefore, in the CMOS device with the conventional poly-SiGe gate, a single poly-SiGe layer is not used, but instead a stacked gate structure having a lower polySiGe layer and an upper poly-Si layer is used. The upper poly-Si layer of the stacked gate structure facilitates the silicide process carried out during the CMOS fabrication to thereby increase the conductivity.\nGenerally, the poly-SiGe layer is formed by a CVD method using a source gas of SiH4 and GeH4. To improve the PDE and the boron penetration characteristics of the poly-SiGe PMOS transistor, it is desirable to form the poly-SiGe layer to have a Ge concentration of at least 20% (see IEEE Electron Device Letters, 10(7), 1998, p.247, by W. C. Lee et. al.). However, when the Ge concentration is more than 30% and is deposited by a chemical vapor deposition (CVD) method, the surface roughness abruptly deteriorates. If the concentration of GeH4 gas is increased to thereby increase the Ge concentration in the poly-SiGe layer, the roughness of the surface deteriorates proportionally to the increase in the GeH4 concentration. Therefore, it is difficult to carry out a subsequent photolithography process, and pits are generated in the substrate during a subsequent etching processes. Accordingly, it is known in the art that the Ge concentration of the poly-SiGe layer for use in the gate is preferably in the range of 20-30%.\nAfter the sequentially stacked poly-SiGe and poly-Si layers are formed as a gate electrode, thermal processes such as gate re-oxidation, silicon nitride deposition, and activation annealing (described above with respect to FIGS. 1A-C) are carried out.\nHowever, at this time Ge diffuses from the lower poly-SiGe layer into the upper poly-Si layer. Accordingly, the Ge concentration of the poly-SiGe layer is reduced from the original Ge concentration (see IEEE Electron Device, 47(4), 2000 p.848 by Y. V. Ponomarev). To account for the Ge diffusion, the Ge concentration in the poly-SiGe layer when it is formed may be set to more than 30%. However, as stated above, this undesirably increases the surface roughness. On the other hand, if the poly-SiGe layer is formed with a Ge concentration of 20-30%, the Ge concentration in the resultant PMOS transistor is reduced to less than 20% after the final CMOS process, thus having little or no effect in controlling the PDE or the boron penetration.\nTherefore, there is a shortcoming in the conventional semiconductor device wherein the gate electrode has the stacked poly-SiGe and poly-Si layers.\nAccordingly, it would be desirable to provide an improved semiconductor device having a stacked poly-SiGe and the poly-Si layer with reduced Ge diffusion. It would further be desirable to provide a method of manufacturing a semiconductor device having a stacked poly-SiGe and the poly-Si layer with reduced Ge diffusion."}
{"text": "1. Field of the Invention\nThe present invention relates to gradation-reduction processing technology for reducing the number of gradations for image data.\n2. Description of the Related Art\nIn recent years, as computer output devices, printers that eject ink from nozzles of a printing head are becoming very popular. For example, as disclosed in Unexamined Patent No. 2002-144551, by using quality-enhancing ink for improving the quality of printed matter, there is an attempt to improve quality of printed matter by improving coloring, moisture resistance, and light resistance, and suppressing glossiness. This kind of quality-enhancing ink is almost transparent, so with the gradation-reduction processing that is performed for generating dot data that shows the formation status of transparent dots on each pixel, even more than data precision, it is highly desirable to have shortened time required for the data generating process.\nHowever, in the past, there had not been consideration of constructing a processing method from the perspective of shortening the time required for gradation-reduction processing that is performed with this kind of dot data generating process for transparent dots."}
{"text": "This invention relates to air metering apparatus and more particularly to such apparatus for use with means for controlling the ratio of air to fuel in a mixture to be combusted in an internal combustion engine.\nThe control of emissions from internal combustion engines and particularly automobile engines has become a major environmental concern. Various federal and state regulatory agencies have promulgated emission standards for certain substances found in the combustion products entering the atmosphere through an engine's exhaust, the most important of these substances being hydrocarbons, carbon monoxide and oxides of nitrogen. To meet emission control standards, various pollution control devices such as catalytic converters and thermal reactors have been developed for use with automobile engines to reduce the quantities of unwanted substances emitted into the atmosphere to within prescribed limits.\nIt has been found that most efficient removal of unwanted substances by pollution control devices is achieved when an engine is operated within a narrow range of air-fuel ratio values for an air-fuel mixture combusted in an engine. Consequently, numerous systems have been developed which attempt to maintain the air-fuel ratio of a mixture to be combusted in an engine within this value range. Examples of systems of this type are disclosed in U.S. Pat. Nos. 3,939,654, 3,946,198, 3,949,551 and 3,963,009. While the systems disclosed in these patents do tend to keep the air-fuel ratio for a mixture to be combusted within the value range where maximum efficiency in removal is obtained, this is usually accomplished only by constantly adjusting the air-fuel ratio. Further, overadjustments frequently occur which then require additional corrections and the systems respond to transitory changes in an engine's operating characteristic to make adjustments when none are actually needed."}
{"text": "There is a need to develop compositions and/or methods that are capable of improving plant yield, for example turfgrass, and reducing or controlling the negative influence of pests and associated damage. To this end, the disclosure provides for formulations comprising fluopyram capable of increasing plant yield and reducing damage from pests.\nFluopyram is a fungicide for treating phytopathogenic diseases. The compound is described in U.S. Pat. No. 7,572,818, which is hereby incorporated by reference in its entirety. Fluopyram is currently being used as a foliar fungicide, and it has also shown to be effective as a seed-applied fungicide/nematicide. However, by way of the current disclosure, it is observed that compositions including fluopyram demonstrate unexpected properties on both plant yield and reduction of plant damage."}
{"text": "1. Field of the Invention\nThe present invention relates to communication device mounting for example a global positioning system (GPS) in a mobile phone or other portable terminal.\n2. Description of the Related Art\nIn a GPS measuring the position of a mobile body utilizing satellites (GPS satellites), the basic function of a GPS receiver is to receive signals from four or more GPS satellites, calculate the position of the receiver from the received signals, and inform it to the user.\nThe GPS receiver demodulates a signal from a GPS satellite to acquire orbital data of the GPS satellite and derives its own three-dimensional position from the orbit and time information of the GPS satellite and delay time of the received signal by simultaneous equations. The reason why four GPS satellites giving received signals are required is that there is error between the time inside the GPS receiver and the time of a satellite and that error must be eliminated.\nIn the case of a GPS receiver for consumer use, a positioning computation is carried out by receiving a spread-spectrum signal wave referred to as the “L1 band” or “C/A” (“coarse acquisition” or “clear and acquisition”) code from a GPS satellite (Navstar).\nThe C/A code is a signal obtained by binary phase shift keying (BPSK) modulating a carrier wave (hereinafter referred to as a “carrier”) having a frequency of 1575.42 MHz by a signal obtained by spreading data of 50 bps by a code of a pseudorandom noise (PN) sequence having a transmission signal rate (chip rate) of 1.023 MHz and a code length of 1023, for example, the Gold code. In this case, since the code length is 1023, the C/A code is formed as a code comprised of the code of the PN sequence repeated using 1023 chips as one cycle (one cycle=1 millisecond (msec)) as shown in FIG. 1A.\nThe code of the PN sequence of this C/A code is different for every GPS satellite, but is configured so that which satellite uses which code of the PN sequence can be detected by a GPS receiver in advance. Further, the navigation message mentioned above enables the GPS receiver to learn from which GPS satellites signals can be received at that position and that point of time. Accordingly, in the case of for example three-dimensional positioning, the GPS receiver receives the waves from four or more GPS satellites which can be acquired at that position and that point of time, despreads the spectrum, and performs the positioning computation to find its own position.\nThen, as shown in FIG. 1B, one bit of the satellite signal data is transmitted as 20 cycles' worth of the code of the PN sequence, that is, 20 milliseconds. Namely, the data transmission rate is 50 bps. 1023 chips of one cycle's worth of the code of the PN sequence are inverted between a time when the bit is “1” and a time when the bit is “0”.\nAs shown in FIG. 1C, in a GPS, one word is formed by 30 bits (600 milliseconds). Further, as shown in FIG. 1D, one sub-frame (6 seconds) is formed by 10 words. As shown in FIG. 1E, the word at the header of one sub-frame has a preamble always regarded as a bit pattern even when data is updated inserted into it. The data is transmitted after this preamble.\nFurther, one frame (30 seconds) is formed by five sub-frames. The navigation message is transmitted in data units of this one frame. First, three sub-frames in this one frame of data form information inherent in the satellite referred to as “ephemeris information”. This information includes parameters for finding the orbit of the satellite and a transmission time of a signal from the satellite.\nAll GPS satellites are provided with atomic clocks and use common time information. The transmission time of a signal from a GPS satellite is a one second unit of the atomic clock. Further, the code of the PN sequence of the GPS satellite is generated as a code in synchronization with the atomic clock.\nThe orbit information of the ephemeris information is updated every several hours, but becomes the same information until that update. By holding the orbit information of the ephemeris information in the memory of this GPS receiver, however, the same information can be precisely used for a few hours. Note that the transmission time of the signal from a GPS satellite is updated every one second.\nThe navigation message of the remaining two sub-frames of one frame of the data is information commonly transmitted from all satellites and referred to as “almanac information”. 25 frames' worth of this almanac information become necessary in order to acquire all information. It is comprised of rough position information of each GPS satellite, information indicating which GPS satellites can be used, etc. This almanac information is updated every several months, but becomes the same information until that update. By holding this almanac information in the memory of the GPS receiver, the same information can be precisely used for several months.\nIn order to receive a GPS satellite signal and obtain the above data, first the carrier is removed, then a code of the PN sequence (hereinafter the code of the PN sequence will be referred to as the “PN code”) the same as the C/A code used in the GPS satellite to be received prepared in the GPS receiver is used to acquire the signal from the GPS satellite and despread the spectrum. When phase synchronization with the C/A code can be established and the spectrum is despread, the bit is detected, and it becomes possible to acquire a navigation message including the time information from the signal from the GPS satellite.\nThe signal from the GPS satellite is acquired by phase retrieval of the C/A code. In this phase retrieval, the correlation between the PN code of the GPS receiver and the PN code of the received signal from the GPS satellite is detected. For example, when the correlation value of the result of the correlation detection is larger than the value set in advance, it is judged that the two are synchronized. When it is judged that synchronization has not been established, any synchronization technique is used to control the phase of the PN code of the GPS receiver to synchronize it with the PN code of the received signal.\nAs mentioned above, a GPS satellite signal is a signal obtained by BPSK modulating a carrier by a signal obtained by spreading data by a spread code. Accordingly, in order for a GPS receiver to receive a GPS satellite signal, it is necessary to establish synchronization of not only the spread code, but also the carrier and the data, but the spread code and the carrier cannot be independently synchronized.\nA GPS receiver generally converts the carrier frequency in the received signal to an intermediate frequency in several MHz and performs the synchronization detection processing mentioned above by that intermediate frequency signal. The carrier in the intermediate frequency signal includes a frequency error mainly due to a Doppler shift in accordance with the velocity of the GPS satellite and a frequency error of a local oscillator generated inside the GPS receiver when converting the received signal to an intermediate frequency signal.\nAccordingly, due to these frequency error factors, the carrier frequency in the intermediate frequency signal is unknown, so a frequency search becomes necessary. Further, a synchronization point (synchronization phase) of the spread code in one cycle depends upon the positional relationship between the GPS receiver and the GPS satellite, so is unknown. Therefore, as mentioned above, some sort of synchronization technique becomes necessary.\nThe GPS receiver uses a synchronization technique employing a frequency search for the carrier and a sliding correlator+delay locked loop (DLL)+Costas loop. This will be explained below.\nThe clock for driving the generator of the PN code of the GPS receiver is generally one obtained by dividing the frequency of the oscillation signal of a reference frequency oscillator provided in the GPS receiver. As this reference frequency oscillator, use is made of a high precision quartz oscillator. From the output of this reference frequency oscillator, a local oscillation signal used for converting the received signal from a GPS satellite to an intermediate frequency signal is generated.\nFIG. 2 is a view for explaining this frequency search. As shown in FIG. 2, when the frequency of the clock signal for driving the generator of the PN code of the GPS receiver is a certain frequency f1, phase retrieval of the PN code is performed, that is, the phase of the PN code is sequentially shifted by one chip, the correlation between the GPS received signal and the PN code at the time of each chip phase is detected, and the peak value of the correlation is detected, whereby it is possible to detect the phase at which synchronization can be established.\nWhen the frequency of the clock signal is f1 and there is no synchronized phase in all of the 1023 chips' worth of the phase retrieval, for example the frequency division ratio with respect to the reference frequency oscillator is changed, the frequency of the drive clock signal is changed to f2, and 1023 chips' worth of phase retrieval is carried out in the same way as above. This is repeated by stepwise changing the frequency of the drive clock signal as shown in FIG. 2. The above operation constitutes the frequency search.\nWhen this frequency search detects the frequency of the drive clock signal regarded to be able to be synchronized, the final phase synchronization of the PN code is carried out at that clock frequency. Due to this, even if there is deviation in the oscillation frequency of the quartz frequency oscillator, it becomes possible to acquire a satellite signal.\nHowever, if the above-mentioned technique is used as the synchronization method, this would be unsuitable for fast synchronization in principle. In an actual receiver, in order to compensate for this, it would be necessary to search for the synchronization point in parallel by forming multiple channels. Further, as described above, if a long time is required for the synchronization of the spread code and carrier, the response of the GPS receiver would become slow. This would be inconvenient in usage. Therefore, for the phase synchronization of the spread code, a technique of performing the code synchronization by a digital matched filter using fast Fourier transform (FFT) processing without using the technique of sliding correlation as mentioned above is realized by improvement of the capability of the hardware such as the digital signal processor (DSP).\nSummarizing the problems to be solved by the invention, in recent years, products combining mobile phones or other networked portable terminals and GPS receivers and services for the same have been put into practical use. As cheaper networks are built in the future, the network hardware and GPS systems will become even more integrated, so it is predicted that network services utilizing positioning will increase.\nThe reference frequency oscillator applied to a GPS system basically has a fixed oscillation frequency, but the reference frequency oscillator used in a mobile phone etc. is configured so as to change in frequency according to a change in the base station or other conditions. Accordingly, if simply configuring a wireless communication device combining a network device and GPS system, it would be necessary to mount two types of reference frequency oscillators. However, this would result in the size of the module becoming larger and would also lead to a cost increase. For such a reason, when configuring a wireless communication device combining a network device and GPS system, it is desirable to make joint use of one reference frequency oscillator, that is, a reference frequency oscillator having a variable frequency as used in mobile phones etc. If simply making joint use of such a variable frequency reference frequency oscillator, however, the following disadvantages would arise.\nAs explained above, the transmission signal referred to as the “L1 band” or “C/A code” of a GPS satellite is a signal obtained by BPSK modulating a carrier of 1575.42 MHz by a signal obtained by directly spreading data of 50 bps by the Gold code having a code length of 1023 and a chip rate of 1.023 MHz. Accordingly, in order to receive a signal from a satellite by a GPS receiver, as mentioned above, it is necessary to establish synchronization between the spread code and the carrier and data, but the synchronization of the spread code and the carrier cannot be independently carried out. The GPS receiver converts the carrier frequency to within several MHz and processes it by the intermediate frequency (IF). However, the carrier at the IF includes an error of the local oscillation frequency mainly generated inside the receiver at the time of the Doppler shift due to the velocity of the satellite and the conversion to the IF, therefore the carrier frequency at the IF is unknown. Further, the phase of the spread code depends upon the positional relationship between the GPS receiver and the satellite, so is unknown. Accordingly, when using a reference frequency oscillator changing in frequency according to the conditions, if the frequency changes, the frequency of the IF carrier also changes. If the frequency of the IF carrier changes, the synchronization between the received signal and the spread code can no longer be held. As a result, it becomes impossible to continuously track the satellite."}
{"text": "An in-line image forming apparatus that includes a plurality of photosensitive drums of respective colors arranged in line and forms a color image by successively superposing toner images of the respective colors on photosensitive drums onto each other is put into practical use, as an image forming apparatus of a plurality of colors or full colors using an electrophotographic method.\nPTL 1 suggests a configuration in which a plurality of photosensitive drums are arranged below an intermediate transfer member and a first transfer surface of the intermediate transfer member is arranged in an inclined manner to decrease the distance of the intermediate transfer member from a first transfer portion to a second transfer portion and hence to decrease a period of time until a transfer material is output from an image forming apparatus. Also, since the intermediate transfer member is arranged in an inclined manner, the dimensions in the lateral direction and height direction of the image forming apparatus are decreased in a balanced manner and hence the space in the image forming apparatus is efficiently used.\nA residual toner that is not second-transferred from the intermediate transfer member to a transfer material and remains on the intermediate transfer member is removed from the intermediate transfer member by a cleaning device. A conventionally known cleaning device is a configuration that scrapes the toner on the intermediate transfer member by a cleaning blade serving as a cleaning member that is pressed to the intermediate transfer member. The scraped toner is removed into a container that is provided in the image forming apparatus.\nPTL 2 discloses a configuration in which an intermediate transfer member is integrally formed with a cleaning device as a transfer unit and the transfer unit is detachably mountable on an apparatus body. Further, PTL 3 discloses a configuration in which a container is attachable to and detachable from a cleaning device and the container is individually detachably mountable on an image forming apparatus."}
{"text": "1. Field of the Invention\nThis invention relates to structure for diverting articles out of a continuous conveying path and, more particularly, to a structure for controllably discharging diverted articles to prevent localized accumulation.\n2. Background Art\nIn conventional glass molding operations, molten glass at approximately 2300.degree.-2600.degree. F. is conveyed continuously in discrete volumes called gobs to mold cavities at a molding station. The size of the gob is dictated by the size of the part, such as a bottle or the like, that is being formed on a particular run.\nTypically, a continuous stream of molten glass flows from a furnace and is separated or cut into the gobs, which are conveyed continuously to the molds. As a general rule, once the furnace is tapped and molten glass flow starts, conveying of gobs from the furnace is carried out continuously through a campaign, which may last from 3-5 years.\nIn the event that a mechanical problem is encountered at the molding station downstream of the furnace or in the event that a mold change is required, it may be necessary to interrupt the molding operation. However, the conveying of the gobs will nonetheless proceed and provision must be made to divert the gobs from their normal route taken to the molding station. Handling the diverted gobs has posed a serious problem to the industry. The high temperature gobs must be accumulated in a form that facilitates recycling and cooled without endangering workmen operating the system.\nHeretofore, two types of structure have been used to accommodate the diverted gobs. One of these is called a tap box. A tap box is essentially a steel container with drilled holes to allow the escape of a cooling fluid which is used to quench the gobs collected in the container. However, the industry is moving away from tap boxes principally because they require the supervision of several workmen to keep up with glass flow and are extremely dangerous to work around. The high temperature gobs, when quenched, may tend to explode into shattered particles generally the size of pea gravel and pose a serious health hazard to those in the vicinity of the tap box. Additionally, the glass gobs accumulate to form a solid mass which is difficult to handle and recycle.\nAs an alternative to the tap box, water filled, vibratory conveyors have been used to intercept, cool and accumulate the diverted gobs. An apparatus suitable for use in a glass molding system is shown in U.S. Pat. No. 4,171,948, to Kraus et al. In Kraus et al a trough contains a predetermined level of cooling fluid and is vibrated to continuously convey material through the fluid between an entry and an outlet end, with the latter residing above the level of the cooling fluid. The gobs entering the fluid explode due to thermal shock upon encountering the fluid.\nOne problem with conventional systems, such as that in Kraus et al, is that the chute which directs the gobs onto the conveyor is fixed so that the gobs are directed onto a single location on the conveying trough. With a steady flow of gobs, such a system is highly acceptable. However, during an upset period, wherein all gobs are diverted out of the main stream, maximum reject rate occurs and several problems may develop.\nFirst, the chutes and conveyor may experience a substantial temperature elevation. The gobs contacting the heated chute remain partially melted and may adhere thereto. Rapid buildup as far back as the molding station may occur and with the resulting heat concentration a fire may erupt and/or dangerous overflow of molten glass may occur.\nAnother problem is that as the temperature of the bath increases, the gobs will not effectively explode and distribute themselves in the trough, which action is normally relied upon by systems designers, resulting in a buildup at the entry portion of the conveying trough. At high flow rates, this condition can deteriorate to the point that some of the gobs may not be immersed and may move off of the trough end in clumps at dangerously high temperatures. Again, a serious threat to the health of workmen is apparent.\nStill further, the unexploded gobs do not convey at the same rate through the fluid as the shattered glass particles. This may result in a backup as far as the chute, with the consequences apparent."}
{"text": "1. Technical Field\nThe present disclosure relates to a check valve, and, more particularly, to a hydraulic check valve having two seats.\n2. Discussion of the Related Art\nCheck valves are mechanical devices used for controlling fluid flow in a system, such as a hydraulic circuit. The fluid may be, for example, a liquid or gas, such as oil or steam. Check valves control flow of the fluid in one direction, and include an inlet opening through which the fluid enters, and an outlet opening through which the fluid exits the check valve. The check valve may be used in a system for hydraulic controls, such as in the aerospace industry.\nCheck valves may operate automatically in a system, whereby the check valve is opened by a minimum upstream pressure caused by the fluid flow. This minimum pressure is referred to as the check valve cracking pressure. A check valve may be designed to open at a predetermined cracking pressure.\nA check valve may include a spring bearing against a poppet, which provides a force on the poppet to maintain the check valve in the closed position. Once the cracking pressure is achieved, the poppet is unseated and opens, further compressing the spring. When the pressure of the flow traveling through and on the check valve decreases, for example, to less than the cracking pressure, the check valve closes to prevent further flow through the valve, and also to prevent backflow, that is, reverse flow back to the fluid source. The check valve may be installed between the fluid source and an actuator driven by the fluid flow to control the flow of the fluid between the source and the actuator.\nThe hydraulic system in which the check valve operates can include a relief valve. The relief valve is an over-pressure pathway designed to relieve excess pressure in the circuit, for example, excess pressure on the actuator, caused by an impact on the actuator or equipment failure. When the relief valve opens, fluid is permitted to flow through the relief valve to, for example, a hydraulic tank, so that excess pressure in the system can be relieved. Like the check valve, a relief valve poppet opens at a predetermined cracking pressure to divert fluid through the over-pressure pathway, and will re-seat upon the pressure reaching the valve's re-seating pressure. The relief valve can include a spring bearing against the poppet, which provides a force on the poppet to maintain the relief valve in the closed position.\nUnwanted fluid leakage may occur in areas where the check valve poppet meets the valve body. In order to prevent this leakage, an O-ring seal can be positioned at the point of leakage. However, there is unwanted friction between the check valve poppet and the O-ring seal.\nThe packing seal introduces additional hysteresis which may cause the poppet to exceed its maximum cracking pressure, and excessive pressure drop at rated flow. Also, the poppet may not reseat causing excessive leakage in the reverse direction.\nTherefore, a need exists for a check valve design that prevents unwanted leakage, and also prevents unwanted friction on the check valve poppet."}
{"text": "A common image coding procedure mainly comprises processes of digital signal transform and quantization. The digital signal transform is performed in a unit of image blocks comprising a plurality of pixels, and selectively conforms to common protocols, e.g., Joint Photographic Experts Group (JPEG), Motion Picture Experts Group (MPEG) 2, MPEG 4 and H.264. The quantization is performed after the digital signal transform to scale down the number of pixels of the image block, so as to reduce a number of bits for representing a single pixel number, and thereby reducing a data amount for representing a single block. By performing the quantization on each of the image blocks, the image block can be divided into different levels. Since the image blocks undergoing the digital signal transform and the quantization generally comprise a large amount of zero quantization coefficients (each of the quantization coefficients corresponds to a single pixel), additional bit numbers are needed for storing the zero quantization coefficients, which is extremely unfavorable for processing the image encoding to a large amount of blocks.\nTherefore, an encoding program called run-length coding (RLC) is applied to encode such blocks comprising a large amount of the zero quantization coefficients, so as to reduce the number of bits and the data amount for storing a single block. Basic principles of the RLC are depicted herein with a binary bitstream. Assume a binary bitstream of “00001100”, it is converted to “402120” with the RLC to represent three RLC strings comprising four successive bits of “0”, two successive bits of “1” and two successive bits of “0”. With principles similar to the above, the image coding utilizing the RLC is capable of reducing the number of bits or the data amount for storing each of the image blocks.\nHowever, the image coding utilizing the RLC also has disadvantages. First, after undergoing the digital signal processing and the quantization, a plurality of quantization coefficients of a single block needs to be linearly scanned one by one for performing the RLC. However, when a plurality of successive zero-runs RLC strings and a plurality of successive nonzero-runs RLC strings are generated during a scanning process for performing RLC on the plurality of quantization coefficients of a single image block, the linearly scanning is rather time-consuming. In addition, when a position of a last non-zero quantization coefficient of a single block is confirmed, it is unnecessary to scan the remainder of zero quantization coefficients for the sake of reducing processing time and increasing efficiency."}
{"text": "Firearms in the modern era conventionally employ a bullet or shell which is positioned within the firearm and has both propellent and the projectile in a single unit. Once operatively positioned within a firearm, a trigger mechanism either in a single action or double action, causes the actuation of a hammer or firing pin which strikes the bullet or shell to cause it to discharge and propel the projectile from the barrel of the firearm toward an intended target.\nBecause of the ease of activating the trigger of many modern firearms, conventionally they are also equipped with a trigger safety mechanism or switch. Such trigger safety switches generally have two primary positions. A first such position is in the safety-on mode, wherein the trigger of the firearm is prevented from actuation which thereby prevents the accidental or unintentional actuation of the trigger and firing of the firearm. A second such position of the safety switch is the safety-off position, wherein the trigger on the firearm is operable by a user to fire the weapon. Some automatic weapons may have a third position configuring the weapon for automatic sequential firing of the weapon. However, the safety on such weapons will still have a safety-on position to prevent firing of the weapon.\nWhile such a safety switch mechanism is in widespread use in firearms and works well to prevent the owner or authorized user of the weapon, it does not prevent unauthorized users from taking and firing the weapon since a simple rotation of the safety will place the firearm in the firing position. As a consequence, for example, children gaining access to a firearm frequently are aware of how to move a safety switch to the safety-off position so that the weapon may be fired. Additionally, where the firearm is stolen or taken from the owner by an individual familiar with trigger safety mechanisms, it may be used against the owner.\nAs a consequence, while conventional trigger safety switches employed with modern firearms work well in many instances, they do not prevent the unauthorized use of the weapon by older children and others coming into possession of it. The consequences of such can be significant harm and loss of life to others encountering the individual who has such unauthorized possession of the firearm.\nThe device and method herein provide for a secondary safety switch which is employable on firearms having a rotating member engaged to the safety which is switchable between a firing and non firing position. The secondary safety switch, however, is in a stealth or hidden position running coaxial to the original safety switch and requires a user key to actuate it between positions. It thereby prevents unauthorized actuation of the trigger of the weapon to which it is engaged.\nThe forgoing examples concerning firearms employing safety switches, and issues regarding such, are intended to be illustrative and not exclusive, and they do not imply any limitations on the invention described and claimed herein. Various other limitations of the related art are known or will become apparent to those skilled in the art upon a reading and understanding of the specification below and the accompanying drawings."}
{"text": "Impellers are widely used in a variety of applications to compress a fluid. For example, impellers are often used in air compressor applications for use in generating compressed air to power pneumatic tools and the like. Alternatively, impellers are used to compress a fluid for use in a pressurized system such as in supplying a pressurized fluid stream for use on a fire truck or pumping station. Further yet, such impellers are commonly used in the design and operation of aircraft engines, whereby a compressed fluid stream is provided via an impeller to propel an airplane in a desired direction. In any of the foregoing applications, it is desirable to provide an impeller capable of operating under varying flow conditions to provide a continuous supply of pressurized fluid, regardless of external forces.\nAs can be appreciated from the foregoing discussion, impellers are operable to compress a fluid stream for use in a plurality of applications. As previously discussed, one such application is an air compressor. Conventional compressors typically include an impeller, a diffuser, and a volute, whereby the diffuser is in fluid communication with both the impeller and the volute and is operable to transfer a compressed air stream from the impeller to the volute for use in an external system. The impeller commonly includes a plurality of blades that are operable to receive and compress an external air stream between a hub of the impeller and a stationary shroud. Specifically, the impeller captures the external air stream at an inducer disposed proximate to a leading edge of each blade such that the captured mass air flow is forced between the hub and the stationary shroud through rotation of the impeller. The inducer is generally operable to capture the external air stream and force it between the hub and the stationary shroud as the impeller is rotated due to the generally curved or arcuate shape of the leading edge of each blade.\nAs can be appreciated, as the air stream travels between each of the blades, the shape of the shroud and hub are such that the air stream is compressed prior to reaching the volute. The compressed air stream is received into the diffuser for distribution to the volute prior to being used by an external application such as a pneumatic tool or a vehicle engine or a fuel cell. The diffuser commonly includes a plurality of stationary vanes which are operable to diffuse the air stream from the impeller in an effort to increase the static pressure of the compressed air. Such increases in static pressure generally increase the pressure of the air stream, thereby providing a desired output of pressurized air from the compressor.\nIn compressor design, it is increasingly important to deliver a constant stream of pressurized air to ensure proper operation of an external device. As can be appreciated, interruption of a compressor can cause external devices, such as pneumatic tools, to seize and abruptly stop working. A common occurrence of such compressor failure is impeller blade fracture or blade cracking due to stresses imparted on the impeller blades through compression of an air stream. Such blade facture or cracking impedes the performance of the impeller as the requisite pressurized air cannot be delivered without first replacing the fractured or cracked impeller blade. Conventional air compressors commonly include an impeller disposed within a sealed housing such that replacement and repair of the impeller commonly requires a significant amount of time and accompanying expense when blade fracture or cracking occurs. As can be appreciated, such repairs can be costly both from the standpoint of requiring a replacement impeller and also from the standpoint that the compressor is unusable until the requisite repairs can be completed.\nTo obviate the need for impeller repair, conventional impeller designs have commonly incorporated impeller blades having an increased thickness to stave off blade cracking and fracture. In most cases, such increases in blade thickness come with an aerodynamic penalty. More particularly, by increasing the thickness of each blade in an effort to improve strength characteristics and limit blade fracture and cracking, aerodynamic performance of the blade is sacrificed as thinner blade profiles typically provide for improved aerodynamic performance and efficiency. In this manner, conventional impellers, and impeller blades, suffer from the disadvantage of sacrificing aerodynamic performance to meet requisite strength characteristics.\nTherefore, an impeller incorporating an airfoil or blade which provides adequate structural support while concurrently providing optimum aerodynamic performance of each airfoil or blade is desirable in the industry."}
{"text": "1. Field of the Invention\nThe present invention relates to an inductor, and more particularly, to a high-current inductor for removing noise that originates in electronic equipment or noise that enters electronic equipment, and a method of manufacturing such an inductor.\n2. Description of the Related Art\nWith recent miniaturization and increases in the level of integration and in the frequency of circuits, there has been an increasing demand for miniature surface-mount inductors.\nAn example of this type of inductor is a multilayer inductor disclosed in Japanese Examined Patent Publication No. 57-39521. This multilayer inductor is manufactured by alternately stacking green sheets of ferrite and conductive patterns for forming a coil, baking them together, and making a multilayer unit that incorporates a coil composed of the coil conductive patterns electrically connected in series.\nA conventional wire-wound inductor may also be manufactured by forming electrodes on a substrate for holding a magnetic core, winding a conductive wire on the magnetic core, and connecting terminals of the conductive wire to the electrodes on the substrate.\nIn the conventional methods of manufacturing the multilayer inductor, however, it is necessary to repeat a printing step many times, the manufacturing process is complicated and the manufacturing cost is high. Although the wire-wound inductor has a relatively simple structure, it is difficult to automate the operations of winding the conductive wire on the magnetic core, processing the terminals of the conductive wire, and soldering the processed terminals of the conductive wire to the electrodes on the substrate for holding the magnetic core. Therefore, these operations are manually performed which increases the manufacturing cost of the wire wound inductor."}
{"text": "The present invention pertains to easels of the type used by draftsmen, artists, store displays, etc.,.\nKnown easels typically are tripods with the legs joined at an apex and with a support for the article in place thereon. Such easels are normally of a single configuration and do not permit use as a drafting table.\nIn the prior art, U.S. Pat. No. 3,512,486 shows an easel having a hinged segment to permit it to be used as an ironing board as well as a drawing board. U.S. Pat. No. 1,286,018 shows an easel with a tiltable platform on which an article may be supported in an upright manner or sloped for drafting. U.S. Pat. No. 3,538,318 discloses an easel having telescopic legs while the easel in U.S. Pat. No. 3,199,825 has collapsible legs for stowage purposes."}
{"text": "1. Field of the Invention\nThe present invention relates to methods for producing heterologous polypeptides in trichothecene-deficient filamentous fungal mutant cells. The present invention also relates to such mutant cells of filamentous fungal cells and methods for obtaining the mutant cells. The present invention also relates to isolated trichodiene synthases and isolated nucleic acid sequences encoding the trichodiene synthases. The present invention also relates to nucleic acid constructs, vectors, and host cells comprising the nucleic acid sequences as well as methods for producing the trichodiene synthases. The present invention further relates to mutants cells comprising a marker-free modification of a gene, and methods for obtaining and using such mutant cells.\n2. Description of the Related Art\nTrichothecenes are sesquiterpene epoxides which are named after the fungus Trichothecium oseum from which the first trichothecene was isolated. The trichothecenes T-2 toxin, diacetoxyscirpenol, and deoxynivalenol are most commonly found in agricultural commodities infected with Fusarium species. Interest in these compounds is due primarily to the discovery that trichothecene contamination of foods and feeds may be detrimental in humans and animals.\nTrichothecenes are produced by a sequence of oxygenations, isomerizations, cyclizations, and esterifications leading from trichodiene, which is produced from the cyclization of trans, trans-farnesyl pyrophosphate by the enzyme trichodiene synthase (Desjardins, Hohn, and McCormick, 1993, Microbiological Reviews 57: 595-604).\nThe trichodiene synthase gene (tri5 or tox5) has been cloned from Fusarium sporotrichioides (Hohn and Beremand, 1989, Gene 79: 131-138); Gibberella pulicaris (Hlohn and Desjardins, 1992, Molecular Plant-Microbe Interactions 5: 249-256); Gibberella zeae (Proctor et al., 1995, Molecular Plant-Microbe Interactions 4: 593-601); Myrothecium roridin (Trapp, et al., 1995, Journal of Cellular Biochemistry Supplement 19B: 154); and Fusarium poae (Fekete et al, 1997, Mycopathologia 138: 91-97).\nTri5- mutants of Gibberella pulicaris (Hohn and Desjardins, 1992, supra) and Gibberella zeae (Proctor et al., 1995, supra) have been generated which do not produce trichothecenes.\nOther genes in the trichothecene biosynthetic pathway have been cloned including the tri3 gene, encoding a 15-O-acetyltransferase, from Fusarium sporotrichioides (McCormick et al., 1996, Applied and Environmental Microbiology 62: 353-359); the tri4 gene, encoding a cytochrome P450 monooxygenase, from Fusarium sporotrichioides (Hohn et al., 1995, Molecular and General Genetics 248: 95-102); the tri6 gene, encoding a zinc finger protein involved in the regulation of trichothecene biosynthesis, from Fusarium sporotrichioides (Proctor et al., 1995, Applied and Environmental Microbiology 61: 1923-1930); the tri11 gene, encoding a cytochrome P450 monooxygenase required for C-15 hydroxylation, from Fusarium sporoirichioides (Alexander et al., 1997, Applied and Environmental Microbiology 64: 221-225); the tri12 gene, which encodes an apparent transport protein involved in trichothecene production, from Fusarium sporotrichioides (Alexander et al., 1997, Cereal Research Communications 25: 347-348); and the tri101 gene, encoding a 3-O-acetyltransferase, from Fusarium graminearum (Kimura et al., 1998, Journal of Biological Chemistry 272: 1654-1661).\nIt is an object of the present invention to provide methods for producing polypeptides in mutant cells."}
{"text": "This invention relates to a process for separating water from, or dewatering, organic solid substrates which contain significant amounts of water. The method has specific relevance to fossilized biomass materials such as peat and coal but also is applicable to other biomass substrates such as agricultural and forestry residues (bagasse, corn stover, sawdust, etc.), groups deliberatey grown for their fuel or chemical value, and sewage sludge. The present invention provides an energy efficient method for removing water from these materials prior to their ultimate use of interim conversion. Since the production of energy or fuel is generally the ultimate use for these substrates, an energy efficient dewatering process could be mandatory for those containing a relatively great amount of water. It is quite likely that the invention concept will be applied to solid/liquid separations other than those described above since nothing in the concept of plug flow contacting restricts its use to water or aqueous fluids. Because of the relatively great amount of energy required to evaporate water, compared for example, to hydrocarbons, dewatering processes benefit most by use of the invention rather than evaporative separation methods.\nWater laden materials can often be dewatered by compression, the applicability of which is dependent upon the solid structure of the material. The predominant effect of compressive loading on water laden substrates of the type considered here is compaction of the solid structure which reduces the void volume available to the water. Material with a rigid structure, such as a particulate mixture with a wide size distribution, is usually difficult to compress. Peat, for example, can be economically dewatered down only to about 75% moisture content by mechanical squeezing.\nA commonly used method of dewatering material such as filter cakes is air displacement. This method is most applicable to materials with an open and uniform pore structure. It is very difficult to displace water from the finer void spaces using air because of the tendency for bypassing through the larger channels.\nNonaqueous liquid displacement of water from a porous material has been used commercially only rarely and then with materials that have high value compared to the biomass and similar type substrates considered here. The main drawback to solvent dewatering is the high cost of solvent recovery and/or loss relative to the product. Our copending patent application, Organic Solid Substrate Dewatering Process Utilizing Primary and Secondary Solvents, Ser. No. 280,264, filed July 2, 1981, now U.S. Pat. No. 4,339,882, typifies interest in liquid dewatering brought about by the high energy cost of conventional dewatering methods.\nLiquid extraction, which is the commonly used and well understood separation process most similar to dewatering by solvent displacement, occurs by bringing about contact between differing liquid phases (one possibly intermixed or associated with a solid) so that a redistribution of the components in the original phases will occur. Design of contacting devices for liquid extraction is based on the definition that an equilibrium distribution, the ultimate product of perfectly intimate contacting, will be achieved from a single perfect contacting stage, (see for example, Treybal, Mass Transfer Operations, Ch. 10--Liquid Extraction). It is well known, however, that for packed columns, such as are commonly used in chromatographic equipment, much better separation can be achieved than would be expected if one assumed the column to be a single equilibrium stage. This is explained in chromatographic theory by showing that the moving zone of contact does not reach equilibrium in the sense generally used in chemical engineering unit operations. This invention demonstrates that a similar phenomena can be used in a liquid displacement dewatering process, especially when used in combination with an energy efficient liquid separation method."}
{"text": "The heart of a computer is an assembly that is referred to as a magnetic disk drive. The magnetic disk drive includes a rotating magnetic disk, write and read heads that are suspended by a suspension arm adjacent to a surface of the rotating magnetic disk and an actuator that swings the suspension arm to place the read and write heads over selected circular tracks on the rotating disk. The read and write heads are directly located on a slider that has an air bearing surface (ABS). The suspension arm biases the slider into contact with the surface of the disk when the disk is not rotating but, when the disk rotates, air is swirled by the rotating disk. When the slider rides on the air bearing, the write and read heads are employed for writing magnetic impressions to and reading magnetic impressions from the rotating disk. The read and write heads are connected to processing circuitry that operates according to a computer program to implement the writing and reading functions.\nThe write head includes a coil layer embedded in first, second and third insulation layers (insulation stack), the insulation stack being sandwiched between first and second pole piece layers. A gap is formed between the first and second pole piece layers by a gap layer at an air bearing surface (ABS) of the write head and the pole piece layers are connected at a back gap. Current conducted to the coil layer induces a magnetic flux in the pole pieces which causes a magnetic field to fringe out at a write gap at the ABS for the purpose of writing the aforementioned magnetic impressions in tracks on the moving media, such as in circular tracks on the aforementioned rotating disk.\nIn recent read head designs a spin valve sensor, also referred to as a giant magnetoresistive (GMR) sensor, has been employed for sensing magnetic fields from the rotating magnetic disk. The sensor includes a nonmagnetic conductive layer, hereinafter referred to as a spacer layer, sandwiched between first and second ferromagnetic layers, hereinafter referred to as a pinned layer and a free layer. First and second leads are connected to the spin valve sensor for conducting a sense current therethrough. The magnetization of the pinned layer is pinned perpendicular to the air bearing surface (ABS) and the magnetic moment of the free layer is located parallel to the ABS, but free to rotate in response to external magnetic fields. The magnetization of the pinned layer is typically pinned by exchange coupling with an antiferromagnetic layer.\nThe thickness of the spacer layer is chosen to be less than the mean free path of conduction electrons through the sensor. With this arrangement, a portion of the conduction electrons is scattered by the interfaces of the spacer layer with each of the pinned and free layers. When the magnetizations of the pinned and free layers are parallel with respect to one another, scattering is minimal and when the magnetizations of the pinned and free layer are antiparallel, scattering is maximized. Changes in scattering alter the resistance of the spin valve sensor in proportion to cos Θ, where Θ is the angle between the magnetizations of the pinned and free layers. In a read mode the resistance of the spin valve sensor changes proportionally to the magnitudes of the magnetic fields from the rotating disk. When a sense current is conducted through the spin valve sensor, resistance changes cause potential changes that are detected and processed as playback signals.\nThe need for ever increased data density is pushing researches to develop data recording systems that can read and record ever smaller bit lengths in order to increase the density of data recorded on a magnetic medium. This has led a push to decrease the gap size of a read head such as a GMR head. However, the amount by which such gap lengths can be decreased has been limited by physical limitations of sensors and also by the limitations of currently available manufacturing methods.\nTherefore, there is a need for magnetoresistive sensor design and method of manufacture that can reduce the read gap of the sensor. Such a design must be manufacturable using currently available manufacturing methods and must not sacrifice sensor reliability and robustness."}
{"text": "This relates in general to automated transaction machines (ATMs). For example, such ATMs may include self-service financial transaction terminals for currency or cash dispensing, accepting, or recycling, which are sometimes referred to as “self-service terminals” or simply “terminals”."}
{"text": "The goal of plant breeding is to combine, in a single variety or hybrid, various desirable traits. For field crops, these traits may include resistance to diseases and insects, resistance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination, stand establishment, growth rate, maturity, plant height and ear height, is important. Traditional plant breeding is an important tool in developing new and improved commercial crops."}
{"text": "Personal computer systems in general, and IBM personal computers in particular, have attained widespread use for providing computer power to many segments of today's society. A personal computer system can usually be defined as a desktop, floor standing, or portable computer that includes a system unit having a system processor, a display monitor, a keyboard, one or more diskette drives, a fixed disk storage, an optional pointing device such as a \"mouse,\" and an optional printer. These systems are designed primarily to give independent computing power to a single user or small group of users and are inexpensively priced for purchase by individuals or businesses. Examples of such personal computer systems are sold under trademarks and include IBM's PERSONAL COMPUTER, PERSONAL COMPUTER XT, PERSONAL COMPUTER AT and IBM's PERSONAL SYSTEM/2 Models 25, 30, 50, 55, 56, 57, 60, 65, 70, 80, 90 and 95 (hereinafter referred to as the IBM products PC, XT, AT, and PS/2, respectively).\nThese systems can be classified into two general families. The first family, usually referred to as Family 1 Models, uses a bus architecture exemplified by the AT computer and other so called IBM compatible machines. The second family, referred to as Family 2 Models, uses IBM's MICRO CHANNEL bus architecture exemplified by IBM's PS/2 models 50 through 95. The bus architectures used in Family 1 and Family 2 Models are well known in the art.\nBeginning with the earliest personal computer system of the Family 1 Models, the IBM PC, and through the current Family 2 Models, the system processor of the personal computer is from the Intel 86 family of microprocessors. The Intel 86 family of processors includes the 8088, 8086, 80286, 80386, and 80486 processors commercially available from Intel Corporation. The architecture of the Intel 86 family of processors provides an upwardly compatible instruction set which assists in preserving software investments from previous processors in the 86 family of processors. This upward compatibility preserves the software application base of the personal computers which use this family of processors. A variety of commonly available and well known software operating systems, such as a DOS or an OS/2 operating system, operate on various members of the Intel 86 family of processors.\nDue to the ever increasing functional requirements in today's electronic systems brought about by the myriad of processors and their capabilities, the number of circuits being integrated into an integrated circuit (IC), as well as the number of IC's being integrated on a circuit board and the number of circuit boards being included in a system, such as a personal computer system, are being increased dramatically. Testing of such a densely packed system becomes more and more difficult. A testing standard presented by IEEE (P1149.1), which is identical to the testing standard adopted by Joint Test Action Group (JTAG) of Europe and North America was developed to address this issue. In these standards, a boundary scan technique or test protocol is presented to serially scan all the devices under test with a test data input (TDI) signal from a JTAG control unit, and the returned scan test data output (TDO) signal is looped back to JTAG control unit. By comparing the predefined serial signal pattern sent out from the JTAG control unit with the returned serial signal pattern, after it has been scanned through the identified devices, the operational readiness of all or certain selected system devices is determined.\nIn a system configured for options such as multiple plug-in cards, concern arises when the JTAG boundary test is considered to be implemented. In such an instance, a serial scan signal is scanned through from card to card via card connectors at card slots or at stations where electronic components or devices, including component boards, are plugged in or inserted and from which, as is usually desirous, such components and devices are readily removable. When the system is not fully populated with such components, devices or cards, a need exists to conduct a serial scan without interruption, regardless of which of the card slots or stations is occupied.\nIn solving this problem, there are many potential different approaches. Cost and reliability are two major deciding factors in selecting an implementation. A reliable and cost effective design, which is not overly complex, is essential. When making a test algorithm feasible, in an otherwise already complex system, it is important not to add more reliability problems to the existing system. To do so would defeat the whole purpose of the initial testing. Also it is desirable to accomplish the testing without the need to include additional signal pins."}
{"text": "1. Field of the Invention\nThis invention relates to electric distribution systems and, more particularly, to an electric panelboard in which bus bars are provided with stab terminals located for easy detachable mounting of one or more circuit breakers.\n2. Description of the Prior Art\nIn the past, circuit breakers for mounting on panelboards have been provided with a neutral wire for separate connection to a neutral bus bar on the panelboard. Such a construction, not only adds to the cost of the panelboard-circuit breaker combination, but also results in wiring inconvenience when the circuit breaker is attached to or dismounted from the panelboard."}
{"text": "1. Field of the Invention\nThe invention relates to the field of sensor resources management and tracking fusion and, in particular, to the tracking of vehicles who's identity becomes ambiguous.\n2. Description of Related Art\nTracking moving ground targets by radar from an aircraft in a battlefield situation is a difficult process. First of all, there may be a large number of moving vehicles in the vicinity of the targets of interest. In addition, the terrain and foliage can intermittently block surveillance. Thus sensor management is critical. In most previous tracking instances, the tracker was data driven. Trackers were at the mercy of the data they ingested. The only way to improve performance was to fine tune prediction models, sensor models, and association algorithms. Such fine-tuning led to improved performance, but only marginally. Potentially, trackers could realize much more significant improvements if they could manage their input data stream.\nThus, it is a primary object of the invention to provide a process for improving the ability to track targets using sensor data.\nIt is another primary object of the invention to provide a process for eliminating ambiguities when tracking vehicles.\nIt is a further object of the invention to provide a process for eliminating ambiguities between targeted vehicles and other vehicles that come within close contact with the targeted vehicle."}
{"text": "Hovercrafts are vehicles that travel across land or water just above a cushion of air. Generally, a skirt extends down from the periphery of the vehicle to the land or water in order to entrap the supporting cushion of air. Depending on the size of the vehicle and/or its intended application, it may be necessary to provide a stability system that reduces the effects of roll and pitch forces. For example, a longitudinally extending keel system can be attached to the underside of the vehicle hull along the vehicle's center line. Ideally, the keel stays in contact with the water's surface to prevent air on one side thereof from moving to the other side thereof. In this way, differential pressure cells are created to provide roll stability. Similarly, a laterally extending keel system can be used to provide pitch stability.\nOne such keel system is utilized by the U.S. Navy in its LCAC (landing craft, air cushion) vehicle. As shown in the plan view of FIG. 1, the LCAC has a segmented longitudinal keel system 20 and a segmented lateral keel system 50 contained within the periphery of a hovercraft's skirt (not shown). Longitudinal keel system 20 extends along the direction of craft motion which is indicated by arrow 100. Lateral keel system 50 is transverse to longitudinal keel system 20 and consists of a plurality of inflated cones 52.\nOne of cones 52 of lateral keel system 50 is shown in a side view in FIG. 2 which is a view taken along line 2--2 of FIG. 1. The side view in FIG. 2 depicts a triangle (of height h and base dimension b) that is essentially positioned parallel to craft motion 100. The portion of cone 52 parallel to the load forces determines the cone's ability to resist the buckling forces associated with craft motion 100. This is best understood by examining the relationship used to determine bending stress .sigma..sub.bending, i.e., the force that causes buckling of a cone 52 subjected to craft motion 100. Bending stress .sigma..sub.bending can be defined mathematically as ##EQU1## where M is the bending moment (i.e., the applied force at the cone's tip or apex times the height h of the cone), c is the distance from the outermost surface of a cone to its central axis (i.e., in this case, one half of the cone's width w as indicated in FIG. 1), and I is the moment of inertia. The moment of inertia I for a triangle is defined as ##EQU2## Since .sigma..sub.bending and I are inversely proportional, the b and h values of each cone 52 must be sufficient to keep .sigma..sub.bending from exceeding a threshold at which the (inflated) cone will buckle.\nThe existing longitudinal keel system is shown in greater detail in isolation in FIG. 3 and in cross-section in FIG. 4 which is taken along line 4--4 of FIG. 3. The longitudinal keel system is attached to the underside of the LCAC's hull which is shown in part and referenced by numeral 10. The longitudinal keel system consists of upper inflatable bladder sections 20A, 20B and 20C (extending fore to aft) attached to the underside of hull 10. Each upper inflatable bladder section 20A, 20B and 20C is trough-shaped and is inflated or pressurized by the craft's lift/fan system (not shown). Attached to each section 20A, 20B and 20C is a corresponding lower trough-shaped inflatable bladder section 22A, 22B and 22C. Inflation of sections 22A, 22B and 22C is accomplished by the provision of a plurality of inflation holes 24 (see FIG. 4) passing between the upper and lower inflatable lower bladder sections 20A and 22A. Thus, inflation of bladder sections 22A, 22B and 22C is accomplished when air passes through these holes. Port and starboard longitudinally extending blades 26 and 28 are attached to the lower portion of sections 22A, 22B and 22C. Blades 26 and 28 are flexible and oppose one another to define a channel therebetween. Holes 30 between blades 26 and 28 in the lower portion of lower inflatable bladder sections 22A, 22B and 22C provide an air exhaust and allow for the removal of water that might seep into the inflated sections while blades 26 and 28 cooperate to prevent water from entering holes 30.\nAlthough accomplishing its intended function of roll stability, the longitudinal keel system shown and described above has been susceptible to environmental abrasion and impact damage that propagates forward and aft from the point of damage. Since the aft-most lower inflatable bladder section 22C is the section most likely to be damaged, such damage tends to propagate forward to section 22B. Damage greatly affects the craft's roll stability and is further costly to repair owing to the location of the keel system."}
{"text": "1. Field of the Invention\nThe present invention relates to sliding roof mechanisms for automobiles, especially passenger cars, and more particularly to a manually operated crank drive for sliding roofs of automobiles which are openable both upwardly, to serve as a roof vent, and retractable under the rearwardly adjacent stationary roof portion, to serve as a sliding roof panel, or so-called sun roof.\n2. Description of the Prior Art\nVarious manually operated two-way openable sliding roofs for automobiles are known from the prior art. Also known is the use of a crank drive for the operation of this type of sliding roof in conjunction with a push-pull type drive cable which is attached to a transverse bridge member connected to the rear edge of the sliding roof panel by means of a lifting linkage, so that the horizontal motion of the bridge member is converted into both a vertical pivoting motion of the sliding roof panel and a retracting motion of the latter, after it has been lowered to a retracting motion of the latter, after it has been lowered to a level below the stationary roof portion to the rear of the sliding roof panel. Such a mechanism is disclosed, for example, in U.S. Pat. No. 3,829,155.\nIn connection with crank drives of the above-mentioned type, it is desirable to be able to swing the crank handle into a scuttled, i.e. partially hidden position, for reasons of convenience and safety of the driver and his passengers, as well as for reasons of appearance and/or in order to block the drive mechanism against movements movements. A special recess is provided for this purpose on the stationary roof structure on which the crank drive is mounted. The crank can be swung into this scuttle recess from below, around a horizontal pivot pin connecting the crank arm to the vertically oriented crank shaft of the drive. The scuttled crank handle locks the crank drive against movements resulting from forces acting on the sliding roof panel, as for example, when wind pressure tends to push the upwardly pivoted panel downwardly, toward its closed position.\nA crank drive of this type has the disadvantage that the crank can be swung into its recess only when it is angularly aligned therewith, which means that the roof panel operating mechanism can be locked in only a limited number of intermediate panel positions, namely in one given position for each crank rotation. This requirement is normally not felt as a disadvantage in connection with sliding roofs which are only openable in one mode, namely by retraction under the stationary roof portion, where a considerable number of crank shaft revolutions are necessary for a movement of the panel between its extreme positions. In the case of a two-way openable sliding roof, however, where the roof panel is also pivotable upwardly from its closed position, the angular restriction on the scuttling of the drive crank represents a considerable shortcoming. Since the lifting motion of the roof panel requires only a comparatively short longitudinal motion of the transverse bridge member, as few as two crank turns will produce the entire opening motion. With known crank drives, the crank has to be moved as much as one-half turn, in order to reach the angular position in which it can be scuttled. This means, for instance, that in the case of two crank revolutions for the lifting movement, the panel can be locked in only one intermediate lifted position. An additional crank revolution from this position either fully opens the roof panel, or brings it to its closed position."}
{"text": "1. Field of the Invention\nThis invention relates to a valve provided for a branch pipe connected to a main pipe, and more particularly to a valve having a function that prevents liquid from stagnating in the branch pipe and the valve when the valve is closed.\n2. Description of the Related Art\nIn order to collect ultra pure water flowing through a main pipe of a ultra pure water production line used, for example, in the semiconductor industry, a branch pipe is connected to the main pipe and a valve is provided for the branch pipe so that ultra pure water is collected from the main pipe by opening the valve.\nFIG. 1 is a view showing an example of a conventional water collecting or sampling (hereinafter referred to as collecting) valve attached to a main pipe.\nIn FIG. 1, reference numeral 20 denotes a main pipe, 21 denotes a branch pipe connected to the main pipe 20, and 22 denotes a valve attached to an end portion of the branch pipe 21. In a conventional collecting valve, when the valve ball 23 in the valve 22 is opened, ultra pure water flowing in the main pipe 20 is introduced into a passageway 25 in the valve case 24 and into a flow passageway 26 provided in the valve 23 to allow to collect ultra pure water from the end opening 27 of the valve case 24.\nIn order to collect ultra pure water flowing in the main pipe 20 to the outside as described above, a valve 22 is provided for the branch pipe 21 connected to the main pipe 20, which valve can be opened or closed.\nHowever, in the case of a valve of the conventional structure, when the valve ball 23 is closed, ultra pure water will stagnate in the branch pipe 21 and passageway 25 in the valve case 24, thereby allowing bacteria to proliferate or small particles to be deposited. Upon resuming the water collection, the stagnant water containing these contaminants will flow out first, and accordingly, there is a problem in that ultra pure water collected must first be blown down and discarded.\nWhat is worse, the stagnation of ultra pure water has an adverse effect upon ultra pure water which flows through the main pipe 20.\nTherefore, the present invention has been made to eliminate the problems of valves of the conventional type described above, and it is an object of the present invention to provide a valve with a function to prevent the stagnation of liquid so as to prevent proliferation of bacteria and deposition of small particles by constructing the valve such that, even when the valve is closed, the liquid does not stagnate in the branch pipe nor in the valve but is allowed to flow continuously therein."}
{"text": "The present invention generally relates to broadcast and multicast services for wireless communication networks, and more particularly, to repair services for broadcast and multicast services.\nIn wireless communication networks, the harsh conditions of the radio communication channel cause errors during transmission of data. Various mechanisms exist to detect and/or correct transmission errors. These techniques include forward error correcting (FEC) codes and automatic repeat request (ARQ) protocols. Both FEC and ARQ operate by adding redundancy to the original information bits to enable detection and/or correction of transmission errors by the mobile station. FEC codes allow the mobile station to correct transmission errors, but the redundant bits reduce the gross bit rate of transmission. An advantage of FEC codes is that no retransmission of information is required to correct errors that occur during transmission unless the error correcting capability of the FEC code is exceeded, and no uplink channel is needed for signaling between the mobile station and the base station. ARQ protocols enable the mobile station to detect frame errors with less redundancy so that a mobile station can request retransmission of incorrectly received frames. Unlike FEC codes, ARQ protocols require an uplink channel for sending retransmission requests. Further, retransmission of missing data consumes bandwidth on the downlink.\nA repair service can be used to provide mobile stations with data that was not correctly received in the original transmission. A repair service is similar to ARQ but uses a dedicated point-to-point (PTP) channel to retransmit incorrectly received data rather than a broadcast channel. A mobile station reports missing segments (e.g., blocks, frames, packets, etc.) to initiate a repair transaction. The mobile station typically sends its repair request after download of the data object is completed. The server then retransmits the missing packets to the mobile station. The repair service uses a PTP uplink and downlink channel for each mobile station. The PTP uplink and downlink channels may be a “shared” channel or “packet based” channel, however, the information transmitted over the PTP channel is specific for each mobile station."}
{"text": "Ribbon bicomponent fibers have conventionally been produced when the fiber is expected to be split into a smaller fiber using mechanical force or by hydroentanglement (e.g., see U.S. Pat. No. 6,627,025 to Yu). The inventors have conceived of using a bicomponent ribbon fiber according to the disclosure provided herein to increase the loft of a nonwoven.\nBulk is often a desirable property for a nonwoven as it transmits a perception of softness and comfort. For example, softness and comfort are desirable for nonwovens used as topsheet or backsheet in diapers. Bulk is also an important characteristic that affects how a nonwoven will absorb, distribute and retain fluids. Good examples are the nonwovens used as an acquisition and distribution layer disposed between the topsheet and the absorbent core of a diaper.\nThe bulk of a nonwoven may be increased with the use crimped fibers in the manufacture of the nonwoven. Traditionally, a nonwoven produced from staple fibers is to use fibers that have been mechanically crimped prior to cutting the fibers to the appropriate length. Such fibers appear to have a zig-zag shape.\nThe typical approach for continuous filaments used in the production of bulky spunbond is to make round fibers using two polymer components having differential shrinkage coefficient when reheated and, to position those fibers in a side-by-side or eccentric way. The differences in shrinkage will force the filament to be twisted in a helix shape. An example of this approach is described in U.S. Pat. No. 5,622,772 to Stoke et al. This approach is sometimes referred to as self-crimped filaments. While this approach can produce a nonwoven with an appearance of high loft, the bulk is easily lost when the nonwoven is compressed by a weight. This is due to the crimps offering little resistance to compression as a result of their shape.\nTherefore there is a need for a self-crimped bicomponent fiber and the spunbond made from such fibers that resist compression, therefore maintaining some of the benefit of the high bulk even when under load."}
{"text": "The present invention relates to a disk array apparatus and a method for controlling the disk array apparatus."}
{"text": "1. Field of the Invention\nThe present invention relates to a focus adjustment apparatus and a focus adjustment method. More specifically, the present invention relates to an automatic focusing technique.\n2. Description of the Related Art\nAccording to a conventional method, an electronic still camera (or video camera) can perform automatic focusing (hereinafter, referred to as “AF”) processing by setting an in-focus position to a lens position where a high frequency component of a luminance signal obtained from an image sensor using a charge coupled device (CCD) can be maximized.\nA conventional scan method includes storing a focus evaluation value based on a high frequency component of a luminance signal obtained from the image sensor while driving the lens in each scanning and designating a lens position corresponding to a maximum value of the stored value as an in-focus position.\nAs another method, a hill-climbing (hereinafter, referred to as “continuous AF”) method includes moving a lens in a predetermined direction so as to increase the focus evaluation value and designating a position where the focus evaluation value can be maximized as an in-focus position.\nAs discussed in Japanese Patent No. 04106485, the time consumed for an AF operation can be reduced by performing the continuous AF processing so as to keep an in-focus state before instructing preparation for a shooting operation, because a moving range of the focusing lens in an AF scan operation to be performed in response to the shooting preparation instruction can be reduced effectively.\nIn the above-described continuous AF processing, if the lens is moved in the predetermined direction to increase the focus evaluation value, an object to be photographed may not be quickly focused in a case where a region to be focused on the image plane is not identified.\nAccording to the method discussed in Japanese Patent No. 04106485, the continuous AF processing is performed in combination with the scan processing to quickly accomplish the focusing operation. However, in this case, a target object to be focused (hereinafter, referred to as a “main object”) is not identified on the image plane to be focused by a user. Therefore, there may be a shooting scene in which the main object cannot be focused."}
{"text": "Data encryption has been employed using a public and private key architecture. In such an architecture, a public key is generated and published, and a private or secret key is also generated and secretly shared with those entities that are allowed to decrypt data. When data is encrypted, it is encrypted using the public key and the only entity that can decrypt the data is one who has the private key.\nAttribute based encryption is a type of encryption that is commonly employed in an environment where a message is encrypted, but the entity that encrypts the message does not necessarily know who will be accessing the data. For instance, in a company or other similar organization, it may be common for a person in the organization to encrypt data so that it can only be viewed by people with sufficient credentials. By way of example, assume that a Vice President of Personnel encrypts a memorandum or other item of information which is only to be viewed by persons on the Vice President's personnel team and the Human Resource Director. Regardless of the specific names of those people, the Vice President may want to encrypt the data so that only that group of individuals (whoever they are), with the appropriate credentials, can view the encrypted information. This is sometimes referred to as credential-based encryption. More generally, the data can be encrypted to any predefined set of attributes. The data is encrypted to a first set of attributes, and the entity that is decrypting the data need only have attributes that are sufficiently close to the first set of attributes. If they are, then the decrypting entity can decrypt the data.\nIn attribute-based encryption, the user's key and ciphertext is labeled with attributes. The user is only allowed to decrypt the ciphertext if the user's key has attributes that sufficiently overlap with attributes on the label of the ciphertext.\nAnother type of attribute-based encryption is referred to as “key policy attribute-based encryption” (KP-ABE). In KP-ABE, each user's private key has an associated access structure which specifies the type of ciphertext the key can decrypt. If the user's access structure is sufficiently close to that specified by the ciphertext, and if the user's key has attributes that are sufficiently close to the attributes on the label of the ciphertext, then the user can decrypt the ciphertext.\nSome work is currently being done in quantum computing. Quantum computing uses the principles of quantum mechanics to represent and manipulate data. Though quantum computers are not yet commercially available, it is believed that quantum computers will be available within the reasonably near future.\nQuantum computers will likely be able to solve many current encryption problems. That is, a quantum computer will likely be able to break a great many current encryptions systems, so that they will no longer be secure.\nThe discussion above is merely provided for general background information and is not intended to be used as an aid in determining the scope of the claimed subject matter."}
{"text": "Solar power or solar energy has been used for many decades for heating dwellings and water and for generating electricity. Because solar power is a renewable energy source much effort has been made to develop systems to use such energy. The costs have been high and the storage of energy has not been very effective. A significant problem is the need for effectively positioning solar concentrators relative to the sun as the sun moves during the day. The best thick plastic-fibers available in the market have been studied but none of these provided the required level of optical transmission in the near-infrared (NIR) and a large fraction of the incident energy is lost in the waveguides after only a few meters propagation. The efficiencies of conventional systems have been low and there is a need for a more efficient and cost-effective system.\nThe method of the present invention provides a solution to the above-outlined problems. More particularly, the present invention relates to a solar power system that has a solar concentrator that has a lens disposed therein. The solar concentrator has a first cone-shape and a bottom. The solar concentrator has a length (l) that is longer than a width (w) at the lens of the solar concentrator. A solid tapering device is disposed at the bottom and has a second cone-shape. The tapering device has a bottom in operative engagement with a first curved glass loop section of a glass rod. The glass rod has a second curved glass loop section and a straight section. The first curved loop section and the second curved glass loop section have a first gap defined therebetween. The second curved glass loop section and the straight glass section have a second gap defined therebetween. The solar concentrator is adapted to receive solar power as rays and conveys and concentrates rays as light towards a focus line or segment at the tapering device. The tapering device is in communication with an upper end of the first curved glass loop section. The first curved glass loop section conveys the light to the second curved glass loop section via and across the first gap. The second curved glass loop section conveys the light to the straight glass section via and across the second gap. While transmitting light, the first curved glass loop section is rotated relative to the second curved glass loop section at the first gap. The second curved glass loop section is rotated relative to the straight glass section at the second gap so that the solar concentrator follows a path of a sun. An outer end of the straight glass section is disposed in proximity to a water surface to heat the water below the water surface.\nIn an alternative embodiment, the outer end is disposed above the water surface.\nIn yet an alternative embodiment, the outer end is disposed above a water surface of water contained in a container and the water surface is maintained between a minimum level and a maximum level.\nIn an alternative embodiment, the light emitted from the outer end converts the water to steam that is conveyed from the container."}
{"text": "Photosensitive elements which can be used in image-reproduction processes are well-known in the graphics arts industry. Such elements are exposed to actinic radiation through an image-bearing transparency, such as a color separation transparency, and developed to produce an image which is either positive or negative with respect to the transparency used. Positive-working elements produce an image which is a duplicate of the transparency through which they are exposed. Negative-working elements produce an image that is the reverse of the transparency through which they are exposed. After imagewise exposure, the photosensitive elements may be developed by washout of soluble image areas, toning tacky image areas with a colorant, peeling apart photoadherent layers, or combinations of these techniques. A series of images may be combined to form a color proof. A useful reference for color proofing methods is Principles of Color Proofing, by Michael H. Bruno, GAMA Communications, Salem,. N.H., 1986.\nIn Cohen and Fan, U.S. Pat. No. 4,282,308, there is described a photosensitive element which is capable of producing reverse, colored images by a dry process without the need for toning with a colorant. The element comprises, in order from top to bottom, a strippable cover sheet, a photoadherent layer containing a colorant, a tacky essentially nonphotosensitive elastomeric layer, and a support. After imagewise exposure to actinic radiation through the cover sheet, the element can be peeled apart by peeling off the cover sheet with the exposed areas of the colored photoadherent layer adherent thereto. The unexposed areas of the colored photoadherent layer remain on the supported elastomeric layer. When the imagewise exposure is through a negative transparency, and the exposed element is peeled apart, a positive image is obtained on the peeled-off cover sheet and a negative image is obtained on the supported elastomeric layer. Imagewise exposure through a positive transparency, followed by peeling apart of the element, yields a negative image on the peeled-off cover sheet and a positive image on the supported elastomeric layer. However, these tacky, essentially nonphotosensitive elastomeric layer containing elements result in high peel force peel-apart systems.\nA need exists for a photosensitive element comprised of layers which have a low peel force and the required adhesion balance so the element can be effectively developed using peel-apart development."}
{"text": "1. Field of the Invention\nThis invention relates to a method of and apparatus for producing hard copy displays of data signals, for example, analog physiological signals, and, in particular, plots of multi-valued functions. One example of a multi-valued function is a circle, the mathematical equation for which is X.sup.2 + Y.sup.2 = r.sup.2 where r is the radius of the circle. This relationship may also be stated as X = .+-. .sqroot.r.sup.2 - Y.sup.2. From this equation it is apparent that for each value of X there exists two values Y and vice versa. There are many, many other multi-valued functions. Curves of multi-valued functions, particularly when plotted on hard copy, are useful in the medical field; for example, such curves can indicate a relationship between heart beat signals derived from several different parts of the body. Such signals from three different planes in the body (referred to as the frontal, horizontal, and sagittal) are useful in the diagnosis of abnormal heart conditions.\n2. Description of the Prior Art\nApparatus for providing single-channel or multi-channel hard copy displays of curves and/or records of the amplitude variations of periodically sampled analog data signals are known in the prior art. In one such known form of apparatus, the values of a plurality of data signals are plotted on a moving photosensitive record sheet by means of a cathode ray tube (CRT) operaating essentially in a single line scan mode. U.S. Pat. No. 3,605,109, assigned to the assignee of the present invention, shows a CRT system of this type. In that system, graphic displays are produced from periodically sampled analog data signals in a so-called connected sampling method the yields substantially continuous, and hence, readily discernible curves, notwithstanding the periodic nature of the sampled data from which the curves are produced. The traces are obtained by producing successive substantially parallel lines on a photosensitive record sheet that is arranged to move past the face of a normally blanked CRT operated in a line scan mode.\nThis prior art apparatus is generally satisfactory. However, it cannot generate a record of a multi-valued function from a series of stored sample data points."}
{"text": "This invention relates to produce labelers and, in particular, to selectively labeling a wide variety of fresh produce by size or grade at high speed.\nWithin the fruit packing industry it has become common practice for fruit to be labeled in large quantities at high speed on sizing/grading equipment. The demand for labeling is driven mainly by the retail industry, requiring product to be identified by variety and/or source and furthermore by size and quality. As a consequence, packers are commonly required to apply a number of different labels to fresh produce during the grading and packing operation. This is generally achieved by installing multiple labeling stations in sequence over the grading conveyors. Two major drawbacks exist with this arrangement: (a) the extended space required to accommodate multiple labeling stations, and (b) the resulting higher costs attached to such installations.\nFor example, if peaches are to be labeled “small, medium or large,” the peaches typically pass through sizing equipment where three banks of labeling equipment are used to apply the appropriate size labels to the sized peaches. The 3 labeling stations (including rotary bellows, etc.) all take up space and are all relatively expensive.\nThere is a clear need for labeling equipment that is more efficient and versatile than the prior art systems that use separate labeling stations for each separate size or grade of produce.\nThe closest prior art known to applicant includes the use of dual cassettes as shown in Rietheimer U.S. Pat. No. 5,645,680 (see FIG. 6). However, Rietheimer has three major weaknesses. First, he uses a complex fixed cam surface housing (as opposed to rotary bellows). Secondly, his system uses conventional knife edge label stripping and is therefore limited to the use of relatively stiff labels. Thirdly, Rietheimer requires that the guide plate of each label cassette discharge each label tangentially to the rotary applicator and parallel to each transfer head as shown in FIG. 6 herein. This geometry limits the number of cassettes usable to two (see FIG. 6). The present invention is capable of using three or more label cassettes without the complex camming mechanism of Rietheimer. The present invention also uses thin and flexible labels, which are usable on more items than stiff labels.\nThe present invention provides, for the first time, a single automatic labeling station capable of applying 3 or more different labels to singulated produce passing through the station. The present invention also provides, for the first time, a single automatic labeling station wherein 3 or more label cassettes interact with a single rotary bellows applicator.\nIn the above example of “small, medium and large” peaches to be labeled, the present invention labels all 3 separate sizes in a single station with a single rotary bellows applicator. The invention reduces most of the prior art machinery required and the space necessary to house the machinery!\nA primary object is to provide an automatic produce labeling apparatus capable of applying 3 or more different labels at a single labeling station and at high speed, i.e., more than 1000 labels per minute.\nA further object is to provide an automatic labeling system which eliminates the need for multiple labeling stations otherwise required by the prior art in applying labels displaying different sizes or grades on produce.\nA further object is to provide labeling apparatus wherein multiple label cassettes interact with a single rotary bellows applicator to apply multiple different labels to produce at speeds in excess of 1000 labels per minute.\nFurther objects and advantages will become apparent from the following description and drawings wherein:"}
{"text": "This invention relates to an airflow distribution system where air is discharged from a relatively thin or shallow plenum chamber into a relatively large work space, chamber, or duct.\nOne application for this invention is in industrial ovens. Typically, there is a thin, shallow plenum on one or more sides of the oven housing, each with an apertured distribution plate directing air into a relatively large, central work chamber where products are supported for operations such as drying, baking, sterilizing, preheating, dehydrating, and curing.\nUniform air distribution within the oven work space is very important in order to maintain uniformity of temperature and processing times throughout. Conventional ovens have had serious drawbacks from the standpoint of uniformity. FIG. 8 is a schematic representation of typical prior art ovens. This is a composite representation showing, on opposite sides, examples of two conventional air distribution plates which have been used prior to the present invention.\nIn FIG. 8, the oven 20 has an inlet 22 with heating coils 24. Air flows to relatively thin, shallow plenum chambers 26 and 28 and passes through distribution plates 30 and 32 into a work chamber 34 from which it exits through outlet 36.\nMost of these prior art ovens have air distribution plates such as that designated 30 for plenum chamber 28 in FIG. 8. It has an array of openings 38. Typically, the air direction and velocity in the shallow plenum chamber 28 carries through the openings 38 into the work space 34 as shown by the arrows 40 in FIG. 8. The angular components of the individual air streams passing through openings 36 are clearly indicated by arrows 40. Another drawback of openings 38 is that they are not adjustable.\nOther prior art ovens have air distribution plates such as that designated 32 shown on the left side of FIG. 8. It has an array of openings 42 each with an angled tab or louver 44. All the tabs or louvers are bent in the same general direction. Airflow is adjustable by bending the tabs and flows into the work chamber generally parallel to the bent tabs in the direction of the arrows 46.\nFor both of the prior art distribution plates (30 or 32) shown in FIG. 8, the airflow into the work chamber is non-uniform causing heating and processing conditions to vary."}
{"text": "The current doctrine of drug delivery requires the solubilization of the drug substance before delivery can occur. In the treatment of disease, the efficacy of many drugs is limited by poor solubility and severe side effects. After decades of research in drug delivery, almost all possible solubilization methods have already been identified. Targeted delivery, i.e., chemically linking a ligand to the delivery vehicle in order to spare normal cells and to kill targeted cells, such as cancer cells and tumors, has been tested. Despite these efforts, it is still very difficult to deliver drugs to specifically targeted sites.\nThe synthesis of ligand-conjugated systems can be difficult. It is even more difficult to achieve targeted drug delivery, while maintaining the solubility of insoluble drugs. One common strategy is to conjugate ligand molecules, which bind to specific receptors on the target cell surface, into the drug delivery system. This often requires the chemical linkage between a ligand and the polymer carrier that encapsulates the anticancer drug. One of the few such systems is based on liposomes, or micelles. However, these bioconjugates polymeric systems are complex and challenging to synthesize. Moreover, it is often difficult to maintain the integrity of the delicate delivery system during storage. Delivery systems are delicate and fragile, particularly during systemic circulation. In addition, every delivery vehicle works on its own, surrounded and solvated by the solvent (i.e., water). As a result, the amount of drug which eventually reaches each targeted cell may be considerably limited. For targeted delivery of bioimaging and radiotherapy agents, it is even more critical for the collective emission of radiation by a cluster of radioisotopes to occur. However, current delivery platforms do not satisfy these requirements.\nThus, there is a need for a delivery system which can target specific disease and/or imaging sites, while still avoiding the problems associated with lack of solubility and efficacy associated with current drug delivery models."}
{"text": "This application claims the benefit of U.S. Provisional Application Ser. No. 60/286,966, filed Apr. 30, 2001."}
{"text": "Database systems that are specifically designed for analyzing very large amounts of data (“analytics workload”) typically do not possess access paths that can address individual data rows efficiently. The only access path offered by most analytical database systems is a table scan. Such systems typically do not have indexes, as database systems targeting transactional workload typical implement them (e.g., B-Trees). This is because the typical analytical workload, at which they excel, hardly requires this and maintaining an index is expensive given the very large data sets which these systems address. As a consequence, retrieving an individual row, e.g., by its key columns, from a table may be as computationally “expensive” as retrieving all rows of the table. This is because there is no index to search for the location of the row in such systems. Hence, such a conventional system compares the search key with the key column of (potentially) all rows in the table.\nAs stated above, this is typically not a problem in analytical workloads that require mostly bulk processing, including when inserting and deleting rows. However, it can become a problem with the increasingly important requirement for real-time analytics. Real-time analytics often requires an analytical database to be synchronized with the operational database very frequently. Approaches such as loading or bulk-synchronizing the analytical database in coarse-grained time intervals typically do not meet the real-time analytics requirements as changes, such as inserted, updated, or deleted rows, in the operational database must find their way into the analytical database very quickly. This may lead to frequent inserts and deletes that only involve a small amount of individual rows. It may be noted that updates can be represented by inserting the modified rows and deleting the old rows. This is particularly a problem for any analytical system that may be updated frequently with a transactional database.\nInserting a row into a table without indexes is typically not a problem, as the new row is simply appended at the end. However, to delete an existing row, its storage location must be known. Without an index the storage location in such conventional systems is determined using a table scan, i.e., by looking at (potentially) all rows. If many rows are deleted from a table at once, the per-row cost may be acceptable. Yet deleting a single row comes at the per-row cost of a full table scan. If done frequently, this is not acceptable for very large tables as typically found in analytical databases.\nIt may also be noted that the same arguments as above apply for hybrid systems which store two copies of the data—e.g., in row-oriented fashion for transactional workloads and in column-oriented fashion for analytical workloads. For such systems, it may also become necessary to synchronize the two copies very frequently. Moreover, in many such systems, one data representation does not possess indexes. Column-oriented representations, for instance, are typically not indexed. In addition to that, it may be noted that both of the described bottlenecks may be independent of the storage representation, i.e., row-vs. column-oriented or other representations."}
{"text": "The Transparent Interconnect of Lots of Links Protocol (Trill) uses an appointed forwarder (AF) mechanism, which is used to appoint a routing bridge (RB) for forwarding packets of a specific virtual local area network (VLAN), so as to prevent local links from forming a loop, because the loop has a very great impact on a network and may cause the network to crash. In the AF mechanism, there is also a VLAN inhibition timer mechanism, which is used to strictly prevent a loop from being formed in a special case or in a case of network status transition.\nThe following describes the AF mechanism and the VLAN inhibition timer mechanism using RB1, RB2, and RB3 on a same link as an example. RB1, RB2, and RB3 are connected through a common switch (bridge). That is, RB1, RB2, and RB3 are connected on the same link, and it is assumed that VLAN 10 is enabled on their ports on this link. RB1 serves as a designated routing bridge (DRB), which is responsible for appointing an AF on the link. RB1 may appoint RB2 as an AF of VLAN 10, and RB1 itself responsible for AF work of another VLAN. In this case, RB2 is responsible for forwarding packets of VLAN 10, and periodically sends a Hello packet to a neighboring RB, indicating that RB2 is the AF of VLAN 10. After receiving the Hello packet, the neighbor of RB2 resets its own inhibition timer for VLAN 10, where the resetting refers to setting the inhibition timer to a default value again. The value of the inhibition timer keeps decrementing with time until it reaches zero. An RB that changes from a non-AF to the AF can start to forward a packet of VLAN 10 only when the inhibition timer for VLAN 10 decrements to 0; otherwise, the RB keeps a state of being inhibited from forwarding packets of VLAN 10. Periodically receiving a Hello packet from an AF of a VLAN prevents from decrementing of an inhibition timer for that VLAN to 0. Even if an RB is appointed as the AF of a VLAN, the RB cannot forward traffic of the VLAN if the inhibition timer for the VLAN does not expire. Therefore, if RB1 and RB3 receive a packet of VLAN 10, because neither RB1 nor RB3 is the AF of VLAN 10, they discard the packet, so as to prevent generation of a loop on the link. After running for a period of time, RB1 sends a Hello packet to appoint RB3 as the AF of VLAN 10, and no longer appoints RB2 as the AF of VLAN 10. After RB3 receives the Hello packet for appointing RB3 as the AF of VLAN 10, if an inhibition timer of RB3 for VLAN 10 does not expire, that is, the inhibition timer for VLAN 10 does not decrement to zero, RB3 still cannot forward a packet of VLAN 10. After receiving the Hello packet that is sent by the DRB and contains AF appointment information, RB2 discovers that RB2 itself is not the AF of VLAN 10. Therefore, RB2 immediately stops forwarding a packet of VLAN 10, and sets, in a Hello packet of VLAN 10 to be sent subsequently, an AF state bit of RB2 for VLAN 10 to zero. In this way, in Hello packets received by RB3 from the neighbors RB1 and RB2, the AF state bits for VLAN 10 are both zero. Because an inhibition timer of a VLAN is reset only when a Hello packet containing an AF state bit being 1 for the VLAN is received, the inhibition timer on RB3 is not reset any longer in this case but keeps decrementing until it reaches zero. When the inhibition timer for VLAN 10 on RB3 decrements to zero, RB3 already appointed as the AF of VLAN 10 starts to forward a packet of VLAN 10. Generally, an inhibition timer is set to be three times as long as a Hello interval, and may be up to 30 seconds.\nHowever, according to the foregoing method, after RB3 receives the Hello packet for being appointed as the AF of VLAN 10, even if RB2 has given up a state of being the AF of VLAN 10 and no longer forwards a packet of VLAN 10, RB3 has to wait for the inhibition timer of VLAN 10 to expire (that is, to decrement to zero) before RB3 can forward a packet of VLAN 10. Therefore, an interruption time of the packet forwarding is relatively long."}
{"text": "Devices such as laser diodes are used in various different applications. Many applications require the laser diode to be on for a relatively short period of time. If a device such as a laser diode is on for a prolonged length of time, due to a fault, there may be a risk that safety is compromised. For example, such devices are often used in consumer products such as mobile phones. Safety may be a particular concern in such devices. Alternatively or additionally the problem may be that if a device is on when it should be off, excess power is consumed. This may be disadvantageous in situations where power is provided, for example, by a battery."}
{"text": "1. Field of the Invention\nThis invention relates to a method for forming a titanium nitride (TiN) barrier layer having a (111) crystallographic orientation on an integrated circuit structure.\n2. Description of the Related Art\nIn the formation of integrated circuit structures, aluminum is used for forming the electrical connections or wiring between active and/or passive devices comprising the integrated circuit structure. Conventionally this involves the use of aluminum which is electrically connected to underlying silicon in the structure. While the aluminum and silicon are electrically connected together, it has become the practice to use intermediate electrically conductive layers interposed between the silicon and aluminum to respectively provide better electrical connection to the silicon, and to provide a physical (metallurgical) barrier between the silicon and aluminum to prevent spiking of the aluminum into the silicon, i.e., migration of aluminum atoms into the underlying silicon, which can interfere with the performance and reliability of the resulting integrated circuit structure.\nConventionally, one method which has been used to accomplish this has been to deposit a layer of titanium over a silicon surface, e.g., to form a with the silicon, and then to anneal the titanium-coated structure in the presence of nitrogen whereby a titanium silicide layer forms over the exposed silicon to form a good electrical contact with the silicon and titanium nitride forms over the titanium silicide as the surface of the titanium layer reacts with the nitrogen atmosphere.\nWhile this method does accomplish the formation of a good electrical contact to the silicon, by formation of the titanium silicide, it often does not result in a satisfactory formation of a barrier layer of titanium nitride over the titanium silicide. This is because the simultaneous formation of both the titanium silicide and the titanium nitride from the same titanium layer results in competing reactions wherein more of the titanium reacts with the silicon, resulting in the formation of a layer of titanium nitride of insufficient thickness to provide the desired barrier protection against aluminum spiking.\nOne prior art solution to this problem has been to form the titanium silicide barrier layer first and then to sputter additional titanium nitride over the titanium silicide or titanium silicide/titanium nitride layer. In this way a sufficient thickness of titanium nitride may be formed to provide the desired barrier layer.\nWhile the above method results in satisfactory formation of a titanium silicide contact layer and a titanium nitride barrier layer over the silicide, and beneath the subsequently deposited aluminum, an problem has been encountered involving electromigration of aluminum atoms in the aluminum layer, during subsequent operation of the integrated circuit structure, if the aluminum layer is not formed with a (111) crystallographic orientation. Such electromigration of the aluminum atoms can result in open circuits in the integrated circuit structure and, therefore, such electromigration must be inhibited or eliminated.\nFormation of titanium nitride by the nitration of titanium will result in formation of a titanium nitride layer having a (111) crystallographic orientation. However, as discussed above, formation of such a titanium nitride from titanium deposited over silicon does not result in formation of a sufficiently thick titanium nitride barrier layer.\nConversely, while sputter deposition of titanium nitride will form the desired thickness of titanium nitride barrier layer, the crystallographic orientation of sputtered titanium nitride is usually either (200) or polycrystalline.\nIt would, therefore, be desirable to form a titanium nitride barrier layer over a silicon surface with a titanium nitride surface having a (111) crystallographic orientation, whereby spiking of aluminum through such a titanium nitride layer would be inhibited or eliminated, yet the formation of aluminum of (111) crystallographic orientation would be promoted by the nucleation sites provided by the (111) crystallographic orientation of the underlying titanium nitride surface."}
{"text": "1. Field of the Invention\nThe present invention relates to diode laser devices that employ laser sources with significant activation thresholds and high temperature sensitivities, which are capable of delivering quasi-continuous or near-continuous wave radiation energy.\n2. Information Disclosure Statement\nDiode light energy sources that have significant activation thresholds and high temperature sensitivities are rarely employed for applications that require laser energy having wavelengths in the visible red to near IR region because the output power at these wavelengths is limited to an ineffective level. Notably, these sources are seldom employed for applications in photodynamic therapy (PDT), which typically require light energy in the 600-760 nm range, or for environmental sensing and monitoring applications, which require light energy in the 3-3.5 xcexcm range.\nPDT is gaining importance both for cancer and non-cancer applications. Typically, this therapy begins with the application of a photosensitizer that may be topically applied, injected, or intravenously introduced to a treatment site. After a suitable time interval depending on certain properties of the photosensitizer, radiation energy in a suitable wavelength band is selectively applied for a predetermined duration and intensity to the target site.\nThe basic concept of PDT is that certain molecules function as photosensitizers that absorb light of certain wavelengths. If light energy of the proper wavelength is delivered to the photosensitizer, it stores the energy from the photons by increasing its energy to a higher level called a triplet state. Inside the body, some of the excited photosensitizer molecules transfer the stored energy to nearby oxygen molecules, exciting them to a higher energy level called a singlet state. Singlet oxygen is a highly reactive gas that rapidly oxidizes essential cellular components that surround it, and in a living cell this oxidation causes necrosis.\nThe choice of wavelength depends on the optical absorption characteristics of the photosensitizer, but red light has been favored for its absorption by the photosensitizer, and ability to reach a large tissue volume. A clinically useful photosensitizer should be non-toxic at useful doses, and should concentrate in diseased tissue by either selective uptake or retention. Also, the photosensitizer should be photochemically efficient and ideally be activated by tissue penetrating light ( greater than 600 nm). Most commonly, photosensitizers used in PDT consist of a hematoporphyrin derivative (HpD) such as porfimer sodium, or Photofrin. Although the mechanism of HpD\"\"s preferential location in malignant cells is uncertain, the total time that the derivatives are retained in the malignant tissue is much longer than in nonmalignant tissue, where it generally clears within 24-72 hours. As a result, there is a xe2x80x9cwindowxe2x80x9d of time in which the physician can exploit the differences in HpD concentrations to cause selective photodegradation of malignant tissue.\nAfter photosensitizer administration, a delay of 24-72 hours allows for HpD to be expelled from healthy tissue, and the malignant tissue is irradiated with visible red light tuned to approximately 630 nm. Shortly after administering treatment, the tumor becomes necrotic and, when effectively treated, the tumor becomes a nonpalpable scab that is usually sloughed off within a few days. A high therapeutic ratio and relative lack of morbidity have made PDT a very attractive form of therapy.\nWhile some forms of disease, such as skin disease, are accessible by a wider variety of light delivery systems, others occur inside the body and light energy must be brought to the treatment site using an optical fiber with a suitably configured tip. PDT is most commonly used to alleviate symptoms and diseases associated with cancers of the brain, head, neck, esophagus, skin, colon, and bladder, although PDT treatment can also be applied to other diseases. An example of the latter is age related macular degeneration, one of the most common causes of blindness in patients over 50 years of age. The xe2x80x9cwetxe2x80x9d form of this disease, responsible for 90% of the severe vision loss associated with this condition, is due to choroidal neovascularization, the proliferation of abnormal blood vessels from the choroid, between the retina and the sclera. Fluid leaking from these blood vessels and development of scar tissue are the main reasons for loss of central vision. PDT can be used to slow down the debilitating effects of the disease. The current treatment is relatively unselective, and the heat used to destroy abnormal vessels can induce severe damage to the overlying retina. Application of PDT for this disease involves an intravenous injection of a photosensitizer such as benzoporphyrin derivative monoacid ring A, followed by photoactivation with non-thermal radiation at a wavelength specific for the dye.\nPDT is advantageous compared to conventional treatments such as surgery because the procedure can often be carried out in a clinical setting, i.e. in treatment rooms rather than operating rooms. Also, PDT is minimally invasive, safe, cost effective, and highly selective.\nHigh laser powers may not be advantageous in PDT because high power densities can cause tissue ablation, and a main mechanism of PDT is to activate the photosensitizer. However, the laser power must be sufficient (3-4 watts) to effectively penetrate tissue and allow treatment distribution over a large area (several cm). Thus, to maximize necrosis, it is advantageous to employ a laser system for PDT that emits radiation having an activating wavelength of a chosen photosensitizer, and is powerful enough to penetrate to a desired depth without damaging the surrounding areas.\nCertain applications, such as the prevention of cell proliferation, that cause restenosis following balloon angioplasty, require delivery of high energy levels through very thin optical fibers (preferably 200 xcexcm diameter). Output powers of 10 to 15 W would be advantageously used, without the risk of thermal damage, because the output radiation is emitted from a long length (several cm) cylindrical tip, that can also be cooled. Short irradiation times are of advantage, as this limits the duration of the treatment.\nContinuous wave (cw) sources are preferred in PDT because they can be evenly applied at a treatment site to yield reproducible results and limited treatment times. Moreover, a cw source maintains the greatest amount of photosensitizer in an excited state that leads to elevated levels of oxygen singlets within the diseased tissue.\nDue to the inefficiency of semiconductor materials, relatively few diode laser sources are presently available for PDT. The heat generated by a diode laser source often limits the power available at a treatment site to an ineffective level. Diode lasers, however, are very effective at longer operating wavelengths such as 808 nm, and 980 nm, because material science has increased the efficiency of the diodes at these wavelengths. This efficiency is lacking at lower wavelengths, i.e. below about 760 nm, and currently, no diode laser source operating at wavelengths used in PDT is able to effectively emit important photosensitizer activating wavelengths, such as 630 nm, 652 nm, or 732 nm, with sufficient average output power of approximately 3 to 4 watts. In fact, for use in PDT, the Food and Drug Administration (FDA) has only approved a KTP frequency doubled Nd:YAG pumped dye laser and an argon-ion pumped dye laser. These systems are inherently difficult to use. They are bulky, can contain toxic dyes, consume large amounts of electricity, need regular maintenance, and cost well over $100,000.\nDiode lasers have distinct advantages over these classical lasers and other solid state lasers including compact size, lower power consumption, and lower maintenance costs. They are much more user friendly, and since a diode laser system typically costs much less than state of the art solid state lasers, laser applications are more available to the public. To take full advantage of the diode laser, there is a need for a diode based laser system capable of transmission of wavelengths within the visible red to near IR region and within the 3.0-3.5 xcexcm range. Thus, the current deficiencies of material science may be overcome and diode lasers could replace solid state lasers in the many applications which require these wavelengths of light energy.\nIt is therefore an object of the present invention to address the need for a diode laser system that may effectively be employed for use in applications that require light energy having a wavelength in the visible red to near IR range or the 3.0-3.5 xcexcm range.\nAnother object of the present invention is to provide a diode laser system that multiplexes outputs from diode laser emitters or emitter groups that are injected with drive currents by transistors.\nStill another object of the present invention is to provide a diode laser system that can emit radiation having a wavelength below 760 nm with sufficient power necessary for activation of common photosensitizers used in PDT.\nIt is an aim of the present invention to provide a pulsed laser device utilizing a diode laser source that simulates quasi-continuous wave output to allow for maximum activation of a chosen photosensitizer, leading to necrosis of diseased tissue.\nIt is a further aim of the present invention to provide a diode laser system that has a narrow and flexible delivery component with a small cross sectional diameter. Maneuverability of the delivery component is thus increased, and it is more easily guided to a treatment site within the mammalian body.\nBriefly stated, the present invention provides a multiplexed medical diode laser system that increases power available for applications that require light energy having a wavelength in the visible red to near IR range or 3.0-3.5 xcexcm range. Radiation from a diode laser emitter or emitter group is multiplexed with radiation from other diode laser emitters or emitter groups into a single optics. The diode laser emitters or emitter groups are injected intermittently with current by transistors and thus, a quasi-cw diode laser source may be available. Since the radiation is multiplexed, diode laser emitters or emitter groups can be operated at optimal duty factors and generate less thermal energy, thereby increasing efficiency while maintaining at least near cw output. In a preferred embodiment, large peak and average output powers are obtained with a combination of fast switching transistors that regulate the duty factor of diode laser emitters or emitter groups, and suitable optics that multiplex the outputs.\nThe above, and other objects, features and advantages of the present invention will become apparent from the following detailed description read in conjunction with the accompanying drawings."}
{"text": "1. Field of Invention\nThis invention relates generally to the identification of a novel gene, Wit 3.0, and its use in wound healing. More particularly, the invention relates to methods and compositions of treatment including, but not limited to, adult skin wound healing, soft tissue wound healing and oral mucosa wound healing.\nSpecifically, the present invention identifies and describes Wit 3.0, a gene, differentially expressed in oral mucosa tissues undergoing wound healing.\nThe present invention also provides for methods of administering Wit 3.0 to incisional and/or excisional soft tissue wounds. Delivery of Wit 3.0 is in antisense form, or alternatively, in sense form.\n2. Description of the Related Art\nBy definition, a wound disrupts the normal epithelial barrier, which is the first line of defense against invading microorganisms. Wound healing typically results after surgery or other injuries to the tissue and involves a complex set of cellular and molecular events to restore the structural integrity of the damaged tissue. Following injury or surgical incision, there are three main stages of wound repair. The initial phase occurs during the first 24 to 48 hours after the event and is characterized by blood coagulation, inflammation and initial wound closure. The middle phase occurs during the first week following the event, and is characterized by cell proliferation and formation of granulation tissue. Granulation tissue is described as small, red, grain-like prominences that form on raw surfaces such as that of wounds or ulcers. Granulations are processes of healing. This middle phase is dependent on new blood vessel formation, or angiogenesis. The last phase, occurring several weeks into the healing and repair process, is characterized by connective tissue synthesis and remodeling.\nFibroblasts in resting tissue are quiescent, immobile cells engaged in minimal biosynthetic activities. After wounding, these cell proliferative and migrate into the wound region where they synthesize and contract a new connective tissue matrix, including collagen, which resembles smooth muscle. Once healing is complete, the wound fibroblasts again become quiescent and regress through apoptosis.\nMigration of fibroblasts to the injury site followed by fibroblast proliferation is stimulated by various growth factors (i.e., epidermal growth factor-EGF, fibroblast growth factor-FGF, transforming growth factor, TGF-beta). These factors also stimulate the synthesis and secretion of extracellular matrix components such as collagen and proteoglycans. At the same time, establishment of new blood vessels (angiogenesis) is occurring and stimulated by fibroblast growth factor (FGF) and vascular endothelial growth factor (VEGF).\nAs healing progresses there is a decrease in the number of active fibroblasts, an increase in extracellular matrix components such as collagen and a decrease in concentration of blood vessels. During this process, the region of the scar is immobilized and collagen fibers are deposited in a near-random fashion.\nFor example, healing of a clean, uninfected incision where the edges have been approximated with sutures begins when the edges fill with clotted blood. The blood clot at the surface becomes dehydrated forming a scab and within a day the epidermis thickens. A few days later, granulation tissue begins to occupy the incision space and fills it in a couple of days. In the event of larger incisions or wounds, there is an increase loss of cells and tissue at the injury site and wound contraction is mediated by myofibroblasts.\nTherefore, the failure of these cellular and molecular events produces abnormal scarring. Examples of abnormal scarring include keloid formation, which is an accumulation of larger than normal amounts of collagen at the wound site creating a protruding scar; proud flesh, which is an abnormal increase in the amount of granulation tissue that blocks reepithelialization; contracture formation, which is immobilized tissue resulting in undesirable fixed, rigid scars that can limit normal range of motion; and fibrosis, which is abnormal connective tissue resulting from myofibroblast activity in the extended wound.\nOther factors also undermine normal wound healing repair mechanisms. In particular, systemic disease, medications, and behavioral factors such as smoking and diet can impede the normal wound repair response. For example, diabetes and peripheral vascular disease impair the formation of healthy granulation tissue such as collagen. Medications such as immunosuppressants can inhibit the inflammatory response and delay wound healing. For example, corticosteroids are immunosuppressants that have dual effects: inhibiting lymphocyte function and inhibiting the synthesis of structural skin proteins, such as collagen.\nIt has been widely recognized that fetal skin and oral mucosa leave minimal scarring. The wound created in these tissues commonly exhibits the rapid initial wound closure likely due to the active approximation of the wound margin within the initial phase, or 24-48-hour period. In contrast, adult skin wounds do not close as rapidly as fetal skin and oral mucosa. It has also been postulated that the adult skin wound lacks initial wound closure mechanisms and tends to create excess tissue contraction and/or fibrosis formation in the later healing stages. Active approximation, similar to that of fetal skin and oral mucosa, in the adult can be achieved by placement of sutures. For example, well-sutured wounds have about 70% of the strength of normal skin. Although, minimal scarring has been observed in fetal skin and oral mucosa, to date, the molecular mechanisms involved in the fetal skin and oral mucosa wound closure have not been elucidated.\nOne-third of the elderly are currently edentulous (without teeth) in either one or both jaws. Formed as a result of wound healing, edentulous mucosa is the portion of the oral mucosa that covers the site where the tooth has been removed. The wound healing process during tooth extraction is thought to follow a similar chronological and physiological pattern as that of typical soft-tissue wound healing as previously discussed.\nTo maximize wound healing, growth factors and small peptides that stimulate the proliferation and biosynthetic activity of cells in the wound have to be isolated and characterized. Existing pharmaceutical creme formulations containing several different growth factors including platelet derived growth factor (PDGF) and transforming growth factor beta (TGFβ) do not ameliorate the situation occurring in the edentulous oral mucosa because the cremes originate from more than one cell type\nAnother disadvantage of the existing formulations containing various growth factors is that the particular factor(s), which offer the greatest benefit in wound healing, cannot be specifically determined. Moreover, effects of multiple growth factors may be potentially adverse on distant organs. As a primary example, angiogenic growth factors enhance the proliferation of blood vessels and hence enhance wound repair. However, these same factors can also enhance neovascularization in areas where it is undesirable, such as accelerating the growth of any benign or malignant tumor.\nThus, an improved treatment is composed of an active agent, which is locally expressed at or near the region of the wound. Such an agent enhances wound healing at the local site of injury but is not deleterious to nearby normal healthy tissues.\nAnother improved treatment administers the active agent to minimize undesired effects on non-wound tissues. For example, a method of administration localizes and limits the active agent to restrict release of the factor to the wound site."}
{"text": "A broadcast network uses the Moving Picture Experts Group-2 Transport Stream (MPEG-2 TS) for delivering multimedia content according to the related art. The MPEG-2 TS is used as a representative transmission technique for delivering a bit stream in which a plurality of broadcast programs (a plurality of encoded video bit streams) can be transmitted without having an error. For example, the MPEG-2 TS is appropriate for digital TV broadcasting, or the like, in the multimedia era.\nFIG. 1 illustrates a layer structure for supporting the MPEG-2 TS according to the related art.\nReferring to FIG. 1, the layer for supporting the MPEG-2 TS includes a media coding layer 110, a sync layer 120, a delivery layer 130, a network layer 140, a data link layer 150, and a physical layer 160.\nThe media coding layer 110 and the sync layer 120 are formed of a format which is easily usable as a basic unit of recording or transmission of media data. Further, the delivery layer 130, the network layer 140, the data link layer 150, and the physical layer 160 record a data block (for example, “AU”) configured by the sync layer 120 in a recording medium or configure a multimedia frame for transmission. The configured multimedia frame is delivered to a subscriber's terminal, or the like, through a predetermined network.\nThe sync layer 120 is composed of a fragment block 122 and an access unit 124, and the delivery layer 130 is composed of an MPEG-2 TS/MP4 132, an RTP/HTTP 134, and a UDP/TCP 136.\nHowever, the MPEG-2 TS has several limits in supporting a multimedia service. That is, the MPEG-2 TS has the limitations of unidirectional communication, non-efficiency in transmission due to a fixed frame size, and a generation of an unnecessary overhead while delivering data by using a transmission protocol and an Internet Protocol (IP) specialized to an audio/video.\nThe MPEG newly suggests the MPEG Media Transport (MMT) standard as one of the multimedia transmission techniques for supporting the multimedia service based on the MPEG technique. The MMT standard is suggested by the MPEG in order to overcome the limits of the MPEG-2 TS.\nThe MMT standard may be applied in order to efficiently deliver hybrid contents through heterogeneous networks. Here, the hybrid contents refer to a set of contents having combined multimedia elements, such as a video/audio/application, or the like. The heterogeneous network means a network in which a broadcast network, a communication network, or the like, are mixed.\nIn addition, the MMT standard defines a transmission technique more friendly to the Internet Protocol (IP) which serves as a basic technique for transmitting the multimedia.\nThe MMT standard typically proves to be an efficient MPEG transmission technique in a multimedia service environment, varying based on the IP; continuous research and standardization thereof have progressed.\nThe MMT standard requires a preparation of a method of providing an efficient MPEG transmission technique in a recent multimedia service environment aiming to provide a hybrid network and hybrid contents.\nTherefore, a need exists for a method and apparatus for configuring a control message to be used in a broadcast system that supports an IP-based multimedia service. The above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "A RE-123-based oxide superconductor (REBa2Cu3O7-x, in which RE represents a rare earth element including Y) which has been found recently exhibits superconductivity at a temperature higher than a liquid nitrogen temperature. It is therefore expected as a promising material from a practical standpoint. There is a strong demand to use the linearly-processed RE-123-based oxide superconductor as a conductor for electric power supply.\nAs a method for processing an oxide superconductor into a linear material, the following method is under consideration. A metal material which is highly strong, heat resistive and is easily processed into a linear material is first formed as an elongated tape. An oxide superconductor is then deposited as a thin film on the tape-shaped metal substrate.\nThe oxide superconductor has electric anisotropy such that the flow of electricity is promoted along a and b crystal axes of the crystals itself, but the flow of electricity is impeded along a c crystal axis of the crystals itself. Accordingly, if the oxide superconductor is provided on a substrate, it is necessary to orient the a and b crystal axes so as to promote the flow of electricity and orient the c crystal axis along another direction.\nHowever, the metal substrate itself is an amorphous material or a polycrystalline material which has a crystal structure significantly different from that of the oxide superconductor. It is therefore difficult to form an oxide superconductor film having high crystal orientation on the substrate. A difference in coefficients of thermal expansion and lattice constants between the substrate and the superconductor may cause distortion in the superconductor or peeling-off of the oxide superconductor film from the substrate during a cooling process to a superconducting critical temperature.\nAs an approach to address these problems, an intermediate layer (i.e., a buffer layer) is first provided on the metal substrate and an oxide superconductor film is formed on the intermediate layer. The intermediate layer consists mainly of MgO, YSZ (yttria-stabilized zirconium), SrTiO3 or the like each having a physical characteristic value, such as a coefficient of thermal expansion and a lattice constant, intermediate with respect to the values of the substrate and the superconductor.\nIn the intermediate layer, the c crystal axis of each crystal thereof are oriented at a right angle with respect to a surface of the substrate, but the a crystal axis (or the b crystal axis) of each crystal thereof are not in-plane oriented along the substantially same direction. Accordingly, the oxide superconductive layer to be provided on the surface of the substrate also has the a crystal axis (or the b crystal axis) each of which are not oriented the substantially same direction in-plane direction, thereby failing to increase critical current density Jc.\nThis problem is eliminated by the ion beam assist method (IBAD method). In the IBAD method, constituent particles ejected from a target by sputtering are deposited on a substrate while being irradiated with argon ions, oxygen ions or other ions at the same time emitted from an ion gun at a tilted angle (e.g., 45 degrees). According to the IBAD method, an intermediate layer having high orientation of the c crystal axis and high in-plane orientation of the a crystal axis with respect to a film deposition surface on the substrate can be provided.\nFIGS. 6 and 7 illustrate an example in which a polycrystalline thin film as an intermediate layer is formed on the substrate by the IBAD method. In FIG. 6, a reference numeral 100 denotes a plate-like substrate and 110 denotes a polycrystalline thin film formed on an upper surface of the substrate 100.\nThe polycrystalline thin film 110 is formed by plural fine crystal grains 120, each having a cubic structure, joined together at grain boundaries. The c crystal axis of each crystal grain 120 is oriented at a right angle with respect to the upper surface (i.e., a film deposition surface) of the substrate 100. The a crystal axis and the b crystal axis of each crystal grain 120 are in-plane oriented along the same direction with each other. The c crystal axis of each crystal grain 120 is oriented at a right angle with respect to the film deposition surface (i.e., the upper surface) of the substrate 100. The crystal grains 120 are joined together with the a crystal axis (or the b crystal axis) being angled (i.e., grain boundary angle K illustrated in FIG. 7) at 30 degrees or less.\nAlthough the IBAD method is considered as a highly practical method in that it provides linear materials with excellent mechanical properties and that stable high characteristics are easy to obtain, the intermediate layer formed by the IBAD method (hereinafter, also referred to as a “IBAD intermediate layer”) has been considered not to have high orientation unless the thickness is not less than about 1000 nm. A deposition rate in the IBAD method is as slow as about 3 nm/min since crystal orientation is controlled by ion beam impact on a non-orientated metal tape. This may take a longer time for film deposition and thus is inferior in productivity.\nAs a method for addressing this problem, development and study have been made intensively to use a fluorite-based oxide, such as YSZ and GdZrO (see, for example, Patent Documents 1 and 2) and a rock-salt-based oxide, such as MgO (see, for example, the Patent Document 1).\nIn the former method, because of the simple lamination structure and wider conditions for film deposition, elongation has been achieved earlier. However, since it is necessary to make the intermediate layer thicker, a production rate becomes low. There is also a problem that the internal stress of the film became large which may cause warpage of the substrate.\nThe latter method is expected to fundamentally solve the above-described problem. Since plural extremely thin films having a thickness of not more than several tens of nanometers are laminated together in this method, however, various kinds of technical know-how have been needed in order to keep the same narrow conditions for film deposition over the elongated substrate.    [Patent Document 1] U.S. Pat. No. 6,933,065    [Patent Document 2] International publication No. 2001-040536"}
{"text": "A switching fabric may refer to a network topology in which network nodes interconnect via one or more network switches, such as crossbar switches (hereinafter called “switching elements”). A switching element may connect multiple inputs to multiple outputs in a matrix manner. If the switching element has X inputs and Y outputs, then the switching element may have a matrix with X×Y cross-points, which are places where the “crossbars” on the switching element may cross. There may be a switch at each cross-point. When closed, the switch may connect one of the X inputs to one of the Y outputs."}
{"text": "The present application relates generally to a fixation system for the treatment of the skeletal system. More particularly, the present invention may be applied to treatment of the human spine.\nSpinal fusion is performed to prevent motion between mobile segments of the spine. A variety of reasons exist for performing spinal fusion. The spine may be unstable due to a traumatic injury, surgery, or invasion and destruction of the vertebrae by tumor. Continued motion of particular segments of the spine may cause overgrowth of joint and ligamentous tissue which, in turn, may compress the spinal cord or its nerves. The curvature of the spine may become abnormal and cause deformity or neurological problems. In these instances, it may be desirable to prevent spinal motion at the affected levels.\nThe spine is composed of individual bones, or vertebrae, stacked on top of each other in a column. Each vertebra includes a cylindrical vertebral body, which participates in weight bearing, and an arch of bone (comprising the lamina and spinous process) which protects the spinal cord and its coverings. The bony arch is connected to the vertebral body by two small columns of bone, referred to as the pedicles. The circular canal between the body, the arch, and the pedicles houses the spinal cord and is called the spinal canal. Between adjacent vertebral bodies lie the intervertebral discs. These are cartilaginous structures that function as shock absorbers for the spine. Facet joints connect the bony arches of the spine and permit spinal motion between adjacent vertebrae.\nSpinal instrumentation is employed as an adjunct to successful spinal fusion. The instrumentation immobilizes the spine while the body forms new, solid bone. Spinal fusion usually is performed by surgically exposing the area of the spine to be fused and thereafter preparing the exposed bone by removing soft tissue and ligaments so new bone can form over the area. After the surgical site has been prepared, an autogenic bone graft (from another part of the body, usually the hip) or an allogenic bone graft (from a cadaver) can be implanted in the prepared area so that new bone can form around and within the implant. Implants have been developed in an attempt to avoid the problems associated with acquiring a bone graft implant. Regardless of the type of implant that is used, the chances of achieving a successful fusion are enhanced if motion in the area is minimized or prevented while new bone forms.\nAlthough there have been advances in this area, there remains a need for improved stabilization systems for use in skeletal fixation and boney fusion procedures."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a source of hydraulic fluid under pressure via which hydraulically operated devices such as active suspensions and the like are operated and more specifically to such an arrangement which permits the size of the switching valve associated with a dual pump system to be reduced in size.\n2. Description of the Prior Art\nJP-A-2-123284 discloses a hydraulic fluid supply arrangement which includes two separate pumps, one having a larger displacement capacity than the other. Each of the pumps is arranged to discharge fluid under pressure into separate discharge passages in which respective first and second check valves are disposed. The two discharge passages merge into a single supply passage or conduit at a point downstream of the check valves.\nIn order to control the amount of fluid which is discharged into the supply passage, a drain control valve arrangement is operatively disposed in two drain passages which respectively lead from the first and second discharge passages at locations upstream of the check valves, and is arranged to selectively determine the amount of fluid which is permitted to drain back to a reservoir from which the pumps induct.\nThe drain control valve arrangement can be conditioned to assume a first state (mode 1) wherein the drain passage which is associated with the larger of the two pumps is open and the drain passage which is associated with the smaller of the two, is closed; a second state (mode 2) wherein the drain which is associated with the smaller is open and that associated with the larger pump is closed; and a third state (mode 3) wherein both of the drain passages are closed. Accordingly, in mode 1 only the discharge of the small pump is output while in mode 2 the output of the small pump is drained and replaced with that of the larger one. In mode 3 the combined discharge of both pumps is output (both drains being closed).\nThe drain control valve arrangement can consist of two serially arranged solenoid operated valves. The first and upstream valve is fluidly communicated with the two drain conduits and is such as to select the drain conduit of the larger pump when de-energized and to switch to the drain conduit of the smaller pump when energized. On the other hand, the second and downstream valve is arranged to determine if the drain passage selected by the first upstream valve is to be permitted to communicate with the reservoir or not. In this instance, this second valve is arranged to establish communication between the first valve and the reservoir when energized. In other words, the first valve is used to determine which of modes 1 and 2 is employed while the second valve induces mode 3 when de-energized.\nHowever, this arrangement is such that, when the spool of the second valve arrangement is moved from its communication cut-off position to the one wherein draining is permitted, the flow of hydraulic fluid through the valve produces a force which tends to move the valve spool back toward its communication cut-off position. Accordingly, as the spring is used to bias the valve spool toward its cut-off position, the solenoid must be able to produce a sufficient force to overcome both the spring and the force produced by the flow of hydraulic fluid through the valve in order to be able to move the spool.\nWith the prior art arrangement, depending on the state of the first valve arrangement and the conditions to which the system is subjected, upon the second valve arrangement being switched to a position wherein drainage is permitted, it is possible for the spool of the second valve arrangement can be exposed to the flow produced by the larger of the two pumps. As this relatively large flow produces a force larger than produced that by the output of the smaller of the two pumps, it is necessary to provide a solenoid which can definitely produce a sufficient force to move the spool against the spring force and the relatively large force produced by the large flow to the position wherein the valve is open and drainage is permitted.\nThis of course induces the drawback that this relatively large solenoid, which is capable of overcoming the maximum resistance to spool movement which is apt to be encountered during system operation, increases the overall bulk, weight and cost of the drain control valve.\nA further drawback which has been encountered with the above type of arrangement is that the temperature of the hydraulic fluid is not taken into consideration. As the temperature of the hydraulic fluid varies, the discharge characteristics of the pumps vary along with the control characteristics of the device which is motivated by the supply of hydraulic fluid under pressure.\nAnother problem which has been encountered with the above type of arrangement is that when the pumps are driven by a belt or similar type of connection with the engine crankshaft at high engine speeds, upon the drain control arrangement being conditioned to stop the draining of both pumps, the resistance to rotation which is exhibited by the pumps sometimes leads to a situation wherein driven connection (e.g. belt) slippage occurs."}
{"text": "Within the food packaging industry, just like in other fields, the importance of being able to provide traceability throughout the value chain has recently increased. Here the word traceability refers to the possibility of being able to trace and control for example the origin of the packaging laminate, converting parameters, product process parameters, and parameters for packaging, sterilising and filling in a packaging line. More in detail it refers to the possibility of being able to determine the origin of the base material from which the packaging laminate is made of, from which batch and reel the packaging laminate is originating, from where the product originates (farm/producer, batch), how the product has been treated in the processing plant, and how the filling and packaging has been performed (filling machine type and id, dairy/packaging plant, sterilization parameters, packaging date). The latter information may be used to generate a best-before-date for the package with content.\nTraceability throughout the value chain facilitates quick recall of packages and the possibility of ensuring a high quality of the final product and, in the end, such offers safety for the consumer and consumer trust for the food producer.\nIn practice, traceability is a complex issue. The value chain for e.g. a milk package is long and with several involved parties and numerous process steps. Paper and polymer granules are manufactured and supplied to the packaging laminate producer which in turn, through converting processes, manufactures packaging laminate. The packaging laminate is delivered to the dairy where it is loaded in the filling machines. The dairy also receives milk from one or several nearby farms and processes the milk within the processing plant of the dairy. After processing, the milk is supplied to the filling machine for packaging. The filling machine produces packages and subsequently said packages are distributed to stores where consumers may pick them from the shelves and buy them.\nIt is known within the packaging industry to provide a packaging laminate having magnetisable particles with a possibility for magnetic marking. Some various aspects are disclosed in WO 2006/135313 and EP-0 705 759, in which several applications are disclosed.\nAll relevant traceability data for a packaging material laminate, e.g. batch number and reel number which traces back to base materials and converting processes information, can magnetically be encoded and read if the laminate comprises suitable magnetic particles arranged as strips, marks, fully covered surfaces or by other means. The magnetic mark may function as an information carrier, keeping non-visible information.\nThe latest development of magnetic sensors allows reading without physical contact between the sensor and the layer with magnetic particles. A readings distance up to 1-2 mm is practically possible. This allows the magnetic layer to be located non-visible on the ready made package, e.g. on the opposite side from printing on a paperboard, thus saving space for graphics and written information on the surface of the package."}
{"text": "The present invention relates to an apparatus for generating electric energy by means of using at least one cavity resonator composed of a vibration plate, at least one side-wall unit, a supporting board and at least one piezoelectric substrate.\nAn energy conversion technique from mechanical vibration to electric power is hopeful as one of the techniques for recovering energy to charge. For example, batteries in a sensor system used in a wide range of applications are requested to have a compact and efficient power supply. The energy conversion technique is mainly based upon a piezoelectric effect, a magnetic induction effect or an electrostatic capacitance effect. Thus, a device having almost the same construction as a microphone is conveniently utilized. One of the most popular devices is a floor type, which has piezoelectric substrates mounted on a vibrating board and generates electric power from mechanical vibration by receiving a pressure on the floor. Recently, an energy conversion device with ZnO wires on the nanometer scale, of which each has a piezoelectricity and can be multi-directionally bended in accordance with mechanical vibration of electrode, has been presented (X. Wang et al., Science, vol. 316, pp. 102-105, 2007). However, there exist still many problems in efficiency of the energy conversion, because these conventional techniques make use of nothing but dimensional change of piezoelectric material accompanied by mechanical vibration. In other words, conventional techniques pass over a resonance vibration phenomenon in piezoelectric materials."}
{"text": "The invention refers to a method for operating a voice mail system connected to a telecommunication system, comprising a recorder with a message storage connected thereto.\nIn large telecommunication systems, which may comprise in particular a plurality of terminals or extension, voice mail systems (also called automatic answering machine) are often connected via a private automatic branch exchange to the telephone network and serve for receiving incoming calls when the recipient or party concerned is absent or busy. When the calling party is forwarded to the voice mail system, he is greeted first of all and then asked to leave a message. With most systems, the greeting can be configured individually for each calling party and can in particular consist in mentioning his name. For the storage of the message, after the greeting, which is usually initiated by a suitable signal in the manner of a request, the voice mail system comprises as a rule a recorder with which a message storage is connected in the manner of a mailbox of the respect-tive user. The recipient of the call can, therefore, start the voice mail system later in order to hear the messages stored in the mailbox.\nThe recall of the speech message stored in the voice mail system is, however, too unpracticable for further use and evaluation of the information contained in the speech message, especially in case of commercial use of the voice mail system. The recipient of the speech message is not in a position to edit the information in a suitable manner and to integrate it into his existing information and communication system.\nTherefore, the invention is based on the problem to provide a suitable method for operating a voice mail system, with the help of which the recipient can use in a particularly flexible way a speech message stored in the voice mail system.\nThe problem is solved by the invention by converting a speech message collected by the recorder, storable as a speech information in the message storage by means of a speech-to-text conversion module into a text message."}
{"text": "1. Technical Field\nThe present disclosure relates to a sensor for detecting a target in a non-contact manner and a switch including the sensor.\n2. Description of the Related Art\nConventionally, a sensor is used in an operating switch for a user's operation such as an operation of turning an apparatus on or off. As this kind of sensor, a capacitance-type sensor utilizing electrostatic capacitance is known. For example, Japanese Unexamined Patent Application Publication No. 2009-111996 discloses a capacitance-type contact sensor (touch sensor).\nThe contact sensor disclosed in Japanese Unexamined Patent Application Publication No. 2009-111996 includes an operating switch, which is a first contact sensor unit, a dummy switch, which is a second contact sensor unit disposed close to the operating switch, and a sensor detection unit that detects whether or not effective contact has been made with the operating switch on the basis of output of the dummy switch. In this contact sensor, while the dummy switch is detecting user's contact, output of the operating switch is made ineffective even if the operating switch is detecting user's contact. In this way, only intentional contact is made effective. This makes it possible to prevent the contact sensor from reacting to an operation that is not intended by the user."}
{"text": "1. Field of the Invention\nThe present invention relates to oil collecting devices and more particularly pertains to a new oil collecting and draining device for easy transportation and depositing of oil drained from a vehicle.\n2. Description of the Prior Art\nThe use of oil collecting devices is known in the prior art. More specifically, oil collecting devices heretofore devised and utilized are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nKnown prior art includes U.S. Pat. Nos. 4,802,599; 5,285,824; 956,399; 4,886,233; 5,402,837; and U.S. Des. Pat. No. 351,394.\nWhile these devices fulfill their respective, particular objectives and requirements, the aforementioned patents do not disclose a new oil collecting and draining device. The inventive device includes a housing having a top wall, a bottom wall, a front wall, a back wall, a first side wall, and a second side wall. The top wall is generally divided into a first portion positioned adjacent to the back wall and a second portion positioned adjacent to the front wall. The top wall is collapsed toward the bottom wall and has a lowest point located in a central area of the first portion. The top wall has an aperture therein extending into an interior of the housing. The aperture is located in the central area of the first portion. A pipe is attached to the front wall and is in communication with an interior of the housing. A valve is positioned on the pipe for selectively opening or closing the pipe. A support structure supports the housing above a ground surface.\nIn these respects, the oil collecting and draining device according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in so doing provides an apparatus primarily developed for the purpose of easy transportation and depositing of oil drained from a vehicle."}
{"text": "Hermetically sealed electric terminals provide an airtight electrical terminal for use in conjunction with hermetically sealed devices where leakage into or from such devices, by way of the terminals, is effectively precluded. For hermetically sealed electric terminals to function safely and effectively for its intended purpose, the terminals require that their conductor pins be electrically insulated from and hermetically sealed to the body portion through which they pass and that an optimum air path be established and thereafter maintained between adjacent portions of the pins and opposite sides of the body.\nIn a conventional hermetic terminal assembly, exemplified by U.S. Pat. No. 3,160,460 to Wyzenbeek, a straight, current carrying pin is fixed in place within a lip defining a hole in the terminal body by a fusible glass-to-metal seal. A resilient insulator is bonded to the face of the body beyond the extent of the glass-to-metal seal. The insulator includes outwardly projecting portions bonded to the conductor pins which define a predetermined air path between adjacent portions of the pins and the body member. Such a hermetic terminal construction has been the standard in the industry for four decades.\nThe primary object of the present invention is to provide a hermetic terminal assembly having conductor pins that are rigidly and hermetically secured to the body portion entirely by a resilient plastic which possesses the requisite materials properties, such as dielectric, moisture resistance, resistance to chemical breakdown, to provide for a hermetic seal. In addition to providing a hermetic seal between the conductor pins and the body, the same resilient plastic is bonded to the conductor pins to provide the desired air path between the pins and the face of the body portion.\nAnother object of the present invention is to provide such a terminal that is simple and economical to manufacture, such as by plastic injection molding."}
{"text": "1. Field of the Invention\nThis invention relates to an assembly structured to maintain a secure engagement between a louver and an associated carrier assembly of a vertical blind assembly and includes a support plate connected to and/or defining an upper end of the louver. A retaining clip is removably disposed in clamping engagement with anyone of a plurality of different types of connecting portions associated with the carrier assembly such that attachment between the connecting portion and the support plate is securely maintained.\n2. Description of the Related Art\nUse of vertical blind assemblies to cover windows, sliding doors, etc. has gained wide popularity over the last few years. Conventionally structured vertical blind assemblies normally include a header having an elongated configuration of sufficient length to extend along the upper periphery of the window, door, etc. intended to be covered. The header includes an interior track extending along substantially the entire length thereof and a plurality of carrier assemblies, at least equal in number to the number of louvers, blinds, vanes, etc. associated with the vertical blind assembly. Each of the carrier assemblies include structural features which facilitate easy passage thereof, as well as the louver or blind attached thereto, along the length of the track. A connecting portion is movably connected to each of the carrier assemblies and extends downwardly therefrom to an exterior of the header and in supporting engagement with a corresponding one of louvers.\nThe structural and operative features of each carrier assembly and associated connecting portion is such as to facilitate movement thereof, and the louver supported thereby, along the length of the track of the header assembly, as set forth above. In addition, each of the carrier assemblies are connected to a “tilt rod” also extending along the length of the header assembly and rotatably connected to mechanical linkage of each of the carrier assemblies so as to facilitate selective turning or tilting, in a concurrent, simultaneous fashion, of each of the louvers. Therefore, in typical fashion the plurality of louvers may extend along any portion of the length of the header in overlying relation to the corresponding door, window, etc. and/or be concurrently disposed into any type of slanted or tilted orientation so as to regulate the amount of light passing through the portal. Admittedly, vertical blind assemblies of the type generally described above include distinguishing structural features depending, at least in part, on a particular application for which a vertical blind assembly is intended. In addition, the vertical blind industry is replete with structurally modified components designed to improve the efficiency of the various types of vertical blind assemblies conventionally known and/or commercially available. Efforts to improve vertical blind assemblies include structural variations in the carrier assembly, connecting portion, track configuration, tilt rod, and selectively operable controls for regulating the position and orientation of the plurality of louvers or blinds associated with the vertical blind assembly.\nBy way of example, the structure of the various known carrier assemblies and their corresponding connecting portions may differ significantly. Common to a majority of such carrier assemblies and connecting portions is the ability of the connecting portion to be removably secured to an upper supporting end of a corresponding one of the louvers to which it is intended to be connected. Removable attachment between the connecting portion and the louver is provided to facilitate an effective connection there between while allowing quick and easy removal of individual ones of the louvers for repair or replacement. However, one disadvantage associated with the conventional manner of interconnecting the upper end of the louver to the connecting portion is the tendency for the louvers to become easily detached therefrom. Such undesirable detachment is particularly prevalent when the collection of louvers or blinds are forced along the length of the track and header assembly in a brisk manner, while a user simultaneously attempts to change the tilted or slanted orientation thereof.\nOther factors commonly affecting the stability of the interconnection between the connecting portion and the individual louvers is the weight, configuration and overall structure of the louver. By way of example only, a louver may be formed, at least in part, from heavy material components such as decorative chain links disposed in depending relation from and supported by a support structure disclosed adjacent an upper most end of the louver. Accordingly, when the aforementioned support portion is removably secured to anyone of a plurality of different types of connecting portions, the weight of the louver or blind may be such as to facilitate its detachment from the connection portion. Such inadvertent detachment is also common when the louver is subjected to unusual forces, such as engagement with people or objects, exposure to wind gusts, or rapid and simultaneous movement and tilting of the louvers through operation of the controls of the vertical blind assembly.\nTherefore, there is a significant and long recognized need in the vertical blind industry for an assembly which securely and consistently maintains a supporting interconnection of the connecting portion and individual louvers or blinds associated with the vertical blind assembly. Such an improved retaining assembly should be structured to not interfere with the normal operation of the vertical blind assembly, especially in terms of the collective movement of the louvers or blinds relative to the header assembly and track structure associated therewith. In addition, an improved retainer assembly of the type needed to overcome known disadvantages and problems of the type set forth above should have sufficient structural and operative versatility to be used with any of the different types of connecting structures known or commercially available, especially when such a large number of structurally distinguishable connecting structures are prevalently used in the vertical blind industry."}
{"text": "This invention relates to and is an improvement of the technique described in copending application Ser. No. 848,521, filed Apr. 7, 1986, by which an alphanumeric character is recognized by the geographical features (e.g. bays and lagoons) which comprise that character.\nOptical character recognition (OCR) devices have long been used to sense and interpret alphanumeric characters that, typically, are provided on a printed page. However, as pointed out in the aforementioned application, such devices usually are limited to the extent that characters are recognizable only if they are printed in one (or, at best, a limited few) predetermined font. Printed characters which, nevertheless, may be clear and formed uniformly as, for example, by a typewriter or other printing machine, will not be recognized if they are formed in some other font.\nAnother operating limitation of conventional optical character recognition devices that is mentioned in that patent application relates to what is known as \"line finding\" and \"segmentation\". A line finding operation is carried out by many conventional OCR devices to locate the lines of characters that are printed on a page. This distinguishes character lines from the spaces between lines and usually is implemented by detecting the distribution of pixel information in the vertical direction. A greater concentration of black pixels represents a line of characters and a low concentration represents a space between lines.\nA segmentation operation is intended to locate the spacing between characters in a line of characters. This isolates a block of pixel data, which then is presumed to be limited solely to a character to be identified, whereafter the data block is examined and recognized. Typical segmentation operations are successful if the characters in a line are spaced uniformly and are printed in, for example, Roman-type font. Characters that are printed at angles, such as italics, often cannot be isolated. Rather, the segmentation operation senses portions of adjacent characters, presumes that such portions comprise a single character, and then attempts, unsuccessfully to identify a \"character\" formed of such portions. Similarly, typical segmentation operations often cannot separate (or \"segment\") characters that are smudged or blurred because of the lack of a well-defined space between such characters. Thus, a block of data representing a character to be identified cannot be formed. Likewise, a break that might be present in a character may be interpreted erroneously as a space between adjacent characters, resulting in two separate blocks representing partial characters rather than a single block representing a whole (albeit broken) character.\nThe aforementioned application suggests that these disadvantages of conventional optical character recognition devices are attributed primarily to the fact that, in most such devices, character segmentation (or separation) might not be successful, thus impeding the comparison of a properly scanned character to a reference, or standard geometric form of that character. Significant deviations between the scanned and reference characters, such as differences in font, misalignment of the scanned character, apparent \"connections\" between independent characters, or \"breaks\" in a single character, largely due to improper character segmentation, prevent accurate identification. While various comparison techniques have been proposed heretofore, most optical character recognition methods rely upon a \"template\" comparison of scanned characters in order to identify those characters.\nIn the OCR technique disclosed in the above-mentioned application, predetermined geographical features in a scanned character are detected, and the detected geographical features are compared to a store of particular geographical features of which known alphanumeric characters are comprised. The scanned character is identified as the alphanumeric character having geographical features which compare favorably to the detected geographical features. The relative positions of the detected geographical features with respect to each other is a primary factor in identifying the scanned character.\nThe detected geographical features are lagoons, which are formed as enclosed areas, and bays, which are formed as open areas. The types and relative positions of the bays, as well as the relative positions of those bays to detected lagoons, determine the identity of the scanned character.\nIn the technique of the aforementioned application, the number of \"class 1\" and \"class 3\" points included in the scanned area are detected. A \"class 1\" point is defined as the free end of a link (a link is a portion of a line included in the character), and a \"class 3\" point is defined as a point formed by the juncture of three links. Also considered are \"class 4\" points which are defined as points formed by the juncture of four links; and one \"class 4\" point is equivalent to two \"class 3\" points. The number of \"class 1\" and \"class 3\" points included in the scanned character significantly narrows the set of characters in which the scanned character is included. The geographical features of the scanned character then are compared to the geographical features of the characters in that set. Since geographical features are detected and analyzed, the particular font of which the character is formed is not significant because a given character includes certain minimum geographical features irrespective of its font.\nIn carrying out this geofeature identification technique, pixels are detected as a character is scanned, and the coordinates of each detected pixel (for example, the XY coordinates) are stored. Only the XY coordinate at the start of a link and the XY coordinate at its end need be stored; and these XY coordinates are determined and stored in real time, as the character is scanned. Depending upon the number of links connected to a particular XY coordinate, the class of the point defined by that XY coordinate is established. The presence of lagoons and bays is ascertained by tracing paths, in predetermined directions, from one connected link to the next, that is, from one stored XY coordinate at the beginning of a link to the next stored XY coordinate which, in most cases, is the end point of one link and the beginning point of a connected link. A lagoon is sensed when the traced path returns to a particular XY coordinate included therein; and a bay is ascertained as a function of changes in the X and Y directions as a path is traced."}
{"text": "In recent years there has been a considerable increase in equipment that uses a low supply voltage, such as communication equipment, portable units, etc. This equipment is provided with several internal devices that require circuits suitable for elevating a direct voltage in certain periods of time; among these voltage elevator circuits charge pump devices are widely used.\nA charge pump structure is present in U.S. Pat. No. 5,874,850. The structure comprises generically a single stage circuit or charge section Si or n stages Si connected in series S1 . . . SN by means of the input terminals C1 . . . CN and the output terminals D1 . . . DN as shown in FIG. 1; in this case the respective side terminals A1 . . . AN of the single stages Si are all connected to a first output terminal O1 of an oscillator OSC while the terminals B1 . . . BN are connected to a second output terminal O2 of the oscillator OSC. The input terminal Ci of a single stage is connected to the output terminal Di+1 of the successive stage; the terminal CN of the last stage is connected to the supply voltage VS while the terminal D1 of the first stage, that constitutes the output terminal OUT of the charge pump circuit structure, is connected to a load storage capacitor SC which has the other terminal connected to ground. The signal generated by the oscillator OSC varies between the reference potential GND, in particular the ground, and the potential of the voltage VS. The voltage on the terminal OUT is equal to N+1 times the voltage VS.\nA single stage circuit or charge section Si is shown in FIG. 2; said circuit stage Si comprises a first charge transfer capacitor TC1 and a second charge transfer capacitor TC2 which have first terminals connected to respective side terminals Ai and Bi, two inverters connected together in loop connection so as to form a flip-flop that has inputs Ji, Ki connected respectively to second terminals of the first charge transfer capacitor TC1 and of the second charge transfer capacitor TC2, negative supply terminals connected together with the input supply terminal Ci and positive supply terminals connected together with the output charge terminal Di.\nThe first inverter is constituted of MOS transistors M1 and M2 while the second inverter is constituted of MOS transistors M3 and M4. In addition FIG. 2 shows the bulk diodes Db1, Db2, Db3, Db4 of the respective transistor M1–M4; the diodes are positioned so that the cathodes of the diodes Db1 and Db3 are connected to the terminal Di while the anodes of the diodes Db2 and Db4 are connected to the terminal Ci.\nIf we consider that the charge pump structure is produced with a single stage circuit Si, when a supply voltage VS is applied the capacitor SC is discharged and the voltage in output is initially VS−2VD, where VD is the voltage at the ends of each diode, approximately 0.7V; the capacitor SC is charged through the bulk diodes. When the voltage difference between the output voltage and the supply voltage VS becomes higher than the threshold voltage of the MOS transistors, the same transistors start to conduct and to load the capacitor SC, replacing the bulk diodes in said operation. The capacitors TC1 and TC2 supply the capacitor SC according to whether respectively the first output terminal of OSC is low and the second output terminal of OSC is high or vice versa.\nIt is possible that in off conditions of the charge pump circuit above described it could happen that the difference of the values of the voltages on the terminals of the capacitors TC1 or TC2 associated with the terminals of the respective inverter is very high and leads to the breakdown of the MOS transistors that constitute the inverters. This can occur if for example one of the capacitors TC1 and TC2 discharges much more quickly than the other by means of the leakage currents or if the voltage on the output terminal increases causing the breakdown of the transistors."}
{"text": "1. Field of the Invention\nThe present invention relates to an orthogonal frequency division multiplexed signal transmitting apparatus, an orthogonal frequency division multiplexed signal receiving apparatus and an orthogonal frequency division multiplexed signal transmitting/receiving system, which employ an orthogonal frequency division multiplexing system and also employ the transmitted data having the packet structure to execute the data transmission like a burst and, more particularly, an insertion of a reference signal used to compensate amplitude and phase characteristics of the spatial transmission line.\n2. Description of the Related Art\nIn recent years, the remarkable development of the multimedia network with the Internet taking the key media is attained and information revolutions containing IT (Information Technology) as the central revolution are in progress in all the countries of the world.\nAt present, the multimedia network is advanced mainly by using the computers that are coupled via the cable. In order to achieve the user's convenience, the coupling of the networks via the radio is important. Thus, the products a part of terminals of which is connected by utilizing the radio have been developed, and are available on the market.\nThen, initially the character data are main as the information that are dealt with by the multimedia network, and shortly the image data, the video data, the sound data are also transmitted. As for the video data, the request to transmit the video data having high resolution like the normal TV image is increased.\nIn such circumstances, the realization of the radio network system that has the transmission rate of 20 Mbit/second or more is studied.\nIn this system, the OFDM (Orthogonal Frequency Division Multiplexing) system is employed in the sub-millwave band of 5 GHz, for example, and the packet type suitable for the radio LAN (Local Area Network), i.e., the burst communication system is employed to transmit the information.\nAs the burst communication system, for example, there are IEEE 802.11a (Institute of Electrical and Electronics Engineers 802.11 activities) that is standardized recently, HYPER LAN 2 of ETIS (European Telecommunications Standards Institute), MMAC (Multimedia Mobile Access Communication; wideband mobile access) system that is discussed in the interior, etc. The studies of all systems are put forward to realize the communication system that is suitable for the radio transmission of all multimedia.\nThis OFDM system is the digital modulation system that is employed in the digital terrestrial broadcastings in Japan and Europe, and is known to have the strong transmitting characteristic with respect to the multipath. In addition, it is possible to transmit the information signal having the higher transmission rate by executing the sufficient estimation of the transmission line to compensate the characteristic of the spatial transmission line.\nSuch estimation of the transmission line is carried out, for example, by providing the preamble data period to the head of the data formulated in the packet structure that is transmitted at a burst, then defining a part of the preamble data, which are transmitted together with a plurality of packets, as the symbol assigned to the reference signal, and then executing the compensation of the transmission line characteristic on the receiver side based on the reference signal transmitted in this manner (referred to as the “reference symbol” hereinafter).\nThen, as the second method of executing the compensation of the transmission line characteristic, there is the method of defining the particular sub-carrier of a plurality of sub-carriers, which transmit the data formulated in the packet structure, as the pilot carrier transmitting carrier, then transmitting the data of the packet structure, then examining the characteristic of the transmitted and received pilot carrier on the receiver side, and then compensating of the transmission line characteristic based on the resultant examined characteristic.\nIn addition, as the third method of executing the estimation of the transmission line characteristic, the method of transmitting the reference symbol every predetermined period if the period in which the data are transmitted is long, etc. may be employed.\nFIG. 1 shows major parameters of OFDM data symbols for the multimedia mobile access, in which the data packeted by the OFDM system whose discussion is in progress in IEEE 802.11a, HYPER LAN 2, MMAC, etc., are transmitted like a burst.\nIn the main parameter specifications of these OFDM data symbols, the OFDM signal is generated by using IFFT (Inverse Fast Fourier Transform) whose order is 64.\nThe transmitted information signal that is formulated as the packet structure is supplied to IFFT as the burst signal, the header signal called the preamble being attached to the head of the burst signal, and the reference signal called the reference symbol and used to decode the data being inserted into a part of the preamble signal.\nIFFT generates 53 sub-carriers consequently. Then, 48 sub-carriers of 53 sub-carriers are assigned as the data carrier, and then these data carriers are digital-modulated by BPSK (Bi-Phase Shift Keying), QPSK (Quadrature Phase Shift Keying), 16 QAM (16-Quadrature Amplitude Modulation) or 64 QAM based on the information signal to be transmitted.\nIn the above respective standards, it is decided that the reference signal is inserted in the preamble period and also transmitted in the data period subsequent to the preamble period. That is, particular four sub-carriers of 53 sub-carriers are defined as the pilot carrier (reference carrier), and the decoding reference signal is transmitted by using the pilot carrier. In this case, as shown in FIG. 1, the remaining one sub-carrier is the zero (center) carrier.\nFIG. 2 shows a configuration of a sub-carrier transmitted from the transmitting apparatus.\nAs described above, 4 sub-carriers of 53 sub-carriers are used as the pilot carrier, in other words the decoding reference signal is always transmitted, and the data carrier is decoded based on the transmitting characteristic, which is transmitted by the decoding reference signal, on the receiver side.\nBy the way, IFFT outputs 53 sub-carriers as 64 time-series sampling point signal. At this time, in order to reduce the multipath distortion in the transmission line, the last 8 or 16 sample signal of the 64 sample signal is inserted prior to the 64 sample signal as the guard interval signal.\nFIG. 3 shows the relationship between the signal generated by IFFT and the guard interval.\nMeanwhile, in the transmission of the decoding reference signal by using the reference symbol and the pilot carrier in the apparatus in the prior art, if respective data carriers have the transmission characteristic with the frequency selectivity due to the multipath characteristic of the spatial transmission line, or if the transmitting amplitude characteristic with the frequency selectivity is varied with the time, or if the transmitted burst data length is long, it is difficult to execute the sufficient estimation of the transmission line by using these reference symbols and four pilot carriers. Thus, the false data are caused frequently depending on the characteristics of the transmission line.\nThen, in HiSWAN (High Speed Wireless Access) as one standard specification in the MMAC promoting conference, it is defined that, in the downward communication executed from the master apparatus to the terminal apparatus, the reference symbol signal in addition to the above pilot carrier is applied periodically. Thus, since the transmission line compensation for all carriers can be carried out with respect to the data signals having the long burst period every predetermined interval by employing the inserted reference symbol signal, the good transmission characteristic can be obtained.\nHowever, if all sub-carriers which are transmitted in the particular symbol period are used as the reference signal, the reduction in the symbol period used in the data transmission is caused, which brings about the reduction in the data transmission efficiency with the reduction of the symbol period.\nIn summary, in the transmission of the reference signal in the prior art, the data transmission with the predetermined performance can be carried out in the good spatial transmission line that is not subjected to the influence of the multipath distortion, however the data error rate is increased under the circumstances in which the multipath distortion characteristic of the spatial transmission line is bad, so that it is impossible to assure the sufficient communication path.\nIn such manner, there is the method of inserting the reference symbols periodically to execute the sufficient compensation of the transmission characteristic even in the spatial transmission line that has the large multipath distortion. In this case, the data error rate can be improved, but there is the drawback that reduction in the data transfer rate is brought about."}
{"text": "A sintered magnet (permanent magnet) of a rare earth alloy is generally produced by compacting powder of the rare earth alloy, sintering the resultant powder compact and subjecting the sintered body to aging. At present, two types of magnets, samarium-cobalt magnets and neodymium-iron-boron magnets, are extensively used in various fields. Among others, neodymium-iron-boron magnets (hereinafter, referred to as “R—Fe—B magnets” where R is any of the rare earth elements including Y, Fe is iron and B is boron) are higher in maximum energy product than any of other various types of magnets, and yet relatively inexpensive. Therefore, the R—Fe—B magnets find positive applications to various types of electronic appliances.\nThe R—Fe—B sintered magnet is essentially composed of a major phase of an R2Fe14B compound of tetragonal system, an R-rich phase made of Nd and the like, and a B-rich phase. A transition metal such as Co and Ni may substitute for part of Fe, and carbon (C) may substitute for part of boron (B). R—Fe—B sintered magnets to which the present invention is suitably applied are described in U.S. Pat. Nos. 4,770,723 and No. 4,792,368, for example.\nTo prepare an R—Fe—B alloy from which the magnet described above is produced, an ingot casting process has been conventionally employed. In a normal ingot casting process, rare earth metal, electrolytic iron and a ferroboron alloy as starting materials are subjected to high-frequency melting, and the resultant melt is cooled relatively slowly in a casting mold, to thereby obtain an alloy ingot.\nRecently, a rapid cooling process such as a strip casting process and a centrifugal casting process has attracted much attention in the art. In the rapid cooling process, a molten alloy is brought into contact with a single roll, a twin roll, a rotating disk, the inner wall of a rotating cylindrical casting mold or the like, to permit comparatively rapid cooling of the molten alloy and in this way, prepare a solidified alloy thinner than an alloy ingot from the molten alloy (hereinafter, such a solidified alloy is called “alloy flakes”). The alloy flakes prepared by such a rapid cooling process normally have a thickness in the range of about 0.03 mm to about 10 mm. According to the rapid cooling process, the molten alloy starts to be solidified from a surface thereof brought into contact with the chill roll (roll contact surface), and crystals grow in a columnar shape in the thickness direction from the roll contact surface. The resultant rapidly solidified alloy, prepared by the strip casting process or the like, has a structure essentially composed of an R2Fe14B crystalline phase, having a minor-axis size of about 0.1 μm to about 100 μm and a major-axis size of about 5 μm to about 500 μm, and an R-rich phase dispersed in the grain boundary between the R2Fe14B crystalline phases. The R-rich phase is a non-magnetic phase including a rare earth element R in a relatively high concentration, and has a thickness (corresponding to the width of the grain boundary) of about 10 μm or less.\nCompared to an alloy prepared by the conventional ingot casting process (die casting process) (hereinafter, such an alloy is called an “ingot alloy”), the rapidly solidified alloy has been cooled in a relatively shorter time (cooling rate: 102° C./sec to 104° C./sec). Accordingly, the rapidly solidified alloy has features that the structure is fine and the crystal grain size is small. In addition, the area of the grain boundary is large and the R-rich phase is dispersed broadly in the grain boundary. Therefore, the rapidly solidified alloy has another feature of excelling in the dispersiveness of the R-rich phase. By using the rapidly solidified alloy having these features, a magnet with excellent magnetic properties can be produced.\nAn alloy preparation method called a Ca reduction process (or reduction-diffusion process) is also known. This process includes the following steps. First, metal calcium (Ca) and calcium chloride (CaCl) are added to either a mixed powder including at least one kind of rare earth oxide, iron powder, pure boron powder and at least one kind of ferroboron powder and boron oxide at a predetermined ratio or a mixed powder including alloy powders or mixed oxides of these constituent elements at a predetermined ratio. The resultant mixture is subjected to reduction-diffusion treatment in an inert atmosphere. The resultant reaction product is put into a slurry state, and the slurry is then treated with water, to thereby obtain a solid of an R—Fe—B alloy.\nA block of a solid alloy is herein called an “alloy block”. The “alloy block” will be any of various forms of solid alloys including not only solidified alloys obtained by cooling melts of an alloy ingot prepared by the conventional ingot casting process, alloy flakes prepared by the rapid cooling process such as the strip casting process and the like, but also a solid alloy prepared by the Ca reduction process.\nAn alloy powder to be compacted is obtained by coarsely pulverizing an alloy block in any form by a hydrogen pulverization process, for example, and/or any of various mechanical milling processes (for example, using a disk mill), and finely pulverizing the resultant coarse powder (mean particle size: 10 μm to 500 μm, for example) by a dry milling process using a jet mill, for example.\nThe R—Fe—B alloy powder to be compacted should preferably have a mean particle size of 1.5 μm to 5 μm from the standpoint of the resultant magnetic properties. “The mean particle size” of a powder herein refers to a mass median diameter (MMD) unless otherwise specified. A powder having such a small mean particle size is however poor in flowability and compactibility (including cavity loading capability and compressibility), and thus poor in productivity.\nTo solve the above problem, coating alloy powder particles with a lubricant has been examined. For example, Japanese Laid-Open Patent Publication No. 08-111308 and its corresponding U.S. Pat. No. 5,666,635 (Assignee: Sumitomo Special Metals Co., Ltd.) disclose the following technique. A lubricant of at least one kind of fatty ester in a liquid form is added in an amount of 0.02 mass % to 5.0 mass % to a rough powder of an R—Fe—B alloy having a mean particle size of 10 μm to 500 μm, and the resultant mixture is pulverized with a jet mill using an inert gas, to prepare a fine powder having a mean particle size of 1.5 μm to 5 μm.\nA lubricant improves the flowability and compactibility of powder, and also functions as a binder for imparting rigidity (strength) to a compact. The lubricant however remains in a sintered body as residual carbon causing degradation of magnetic properties, and is therefore required to have good removability. For example, Japanese-Laid Open Patent Publication No. 2000-306753 discloses, as a lubricant having good removability, a depolymerized polymer, a mixture of a depolymerized polymer and a hydrocarbon-based solvent, and a mixture of a depolymerized polymer, a low-viscosity mineral oil and a hydrocarbon-based solvent.\nUse of a lubricant described above contributes to some degree of improvement, but fails in imparting sufficient compactibility. In particular, a powder prepared by the strip casting process, which is not only small in mean particle size but also narrow in particle size distribution, is especially poor in flowability. This causes problems such that the amount of powder loaded in a cavity tends to vary beyond an acceptable range, and that the loading density in the cavity tends to lack uniformity. As a result, the mass and size of the resultant compact may vary beyond an acceptable range, and chips and fractures may be formed in the compact.\nAs another method for improving the flowability and compactibility of a R—Fe—B alloy powder, use of a granulated powder has been attempted.\nFor example, Japanese Laid-Open Patent Publication No. 63-237402 discloses that compactibility can be improved by use of a granulated powder obtained by adding a mixture of a paraffin mixture that is in the liquid state at room temperature and an aliphatic carboxylic acid in an amount of 0.4 to 4.0 mass % with respect to a powder, and kneading and then granulating the resultant mixture. A method using polyvinyl alcohol (PVA) as a granulating agent is also known. Like the lubricant, the granulating agent functions as a binder imparting strength to a compact.\nThe granulating agent disclosed in Japanese Laid-Open Patent Publication No. 63-237402 described above is poor in removability. Therefore, in the case of production of an R—Fe—B sintered magnet, the magnetic properties disadvantageously degrade due to carbon remaining in a sintered body.\nIn the method using PVA, a granulated powder prepared by a spray dryer process using PVA has strong binding force. The resultant granulated powder is too rigid to disintegrate under application of an external magnetic field. Therefore, alloy particles (crystals) fail to be sufficiently aligned in the magnetic field, and as a result, no anisotropic magnet excellent in magnetic properties is obtainable.\nPVA is also poor in removability. Therefore, carbon derived from PVA tends to remain in the resultant magnet, causing degradation in magnetic properties. Debindering may be performed in a hydrogen atmosphere. Even with this treatment, however, it is difficult to sufficiently remove the carbon. Also, due to the excessively strong binding force of PVA, the granulated powder fails to disintegrate under application of a magnetic field and therefore finds difficulty in being aligned.\nAs described above, although various granulating agents have been examined so far, there has not yet been succeeded in development of a granulating agent that has moderate binding force and is excellent in removability. Under this circumstance, a method permitting industrial-scale production of a granulated powder suitably usable for production of an R—Fe—B sintered body has not yet been attained.\nNeeds for smaller/thinner magnets with higher performance have grown. In this situation, development of a method permitting production of small/thin magnets with high performance with high production efficiency is desired. In general, when an R—Fe—B alloy sintered body (or a magnet obtained by magnetizing the alloy sintered body) is machined, the magnetic properties of the machined product degrade due to machining strain. This degradation of the magnetic properties is not negligible for small magnets. In view of this, as the magnet is smaller, it is desired more strongly that a sintered body in a final use shape should be produced with a level of size accuracy high enough to substantially require no machining. Under this circumstance, also, the demand for an R—Fe—B alloy powder material excellent in flowability and compactibility has been increasingly intensified."}
{"text": "1. Field of the Invention\nThis invention relates to InGaN/GaN light emitting diodes (LEDs) and laser diodes (LDs), and more particularly to nonpolar III-nitride LEDs in which the wavelength can be controlled by selecting the barrier thickness of multiple quantum wells (MQWs).\n2. Description of the Related Art\nThe usefulness of gallium nitride (GaN) and its ternary and quaternary compounds incorporating aluminum and indium (AlGaN, InGaN, AlInGaN) has been well established for fabrication of visible and ultraviolet optoelectronic devices and high-power electronic devices. These compounds are referred to herein as Group III nitrides, or III-nitrides, or just nitrides, or by the nomenclature (Al, B, Ga, In)N. Devices made from these compounds are typically grown epitaxially using growth techniques including molecular beam epitaxy (MBE), metalorganic chemical vapor deposition (MOCVD), and hydride vapor phase epitaxy (HVPE).\nCurrent nitride technology for electronic and optoelectronic devices employs nitride films grown along the polar c-direction. However, conventional c-plane quantum well structures in III-nitride based optoelectronic and electronic devices suffer from the undesirable quantum-confined Stark effect (QCSE), due to the existence of strong piezoelectric and spontaneous polarizations. The strong built-in electric fields along the c-direction cause spatial separation of electrons and holes that in turn give rise to restricted carrier recombination efficiency, reduced oscillator strength, and red-shifted emission.\nOne approach to eliminating the spontaneous and piezoelectric polarization effects in Group-III nitride optoelectronic devices is to grow the devices on nonpolar planes of the crystal. For example, in GaN crystals, such planes contain equal numbers of Ga and N atoms and are charge-neutral. Furthermore, subsequent nonpolar layers are equivalent to one another so the bulk crystal will not be polarized along the growth direction. Two such families of symmetry-equivalent nonpolar planes in GaN are the {11-20} family, known collectively as a-planes, and the {10-10} family, known collectively as m-planes.\nThe other cause of polarization is piezoelectric polarization. This occurs when the material experiences a compressive or tensile strain, as can occur when (Al, B, Ga, In)N layers of dissimilar composition (and therefore different lattice constants) are grown in a nitride heterostructure. For example, a thin AlGaN layer on a GaN template will have in-plane tensile strain, and a thin InGaN layer on a GaN template will have in-plane compressive strain, both due to lattice matching to the GaN. Therefore, for an InGaN quantum well on GaN, the piezoelectric polarization will point in the opposite direction than that of the spontaneous polarization of the InGaN and GaN. For an AlGaN layer lattice matched to GaN, the piezoelectric polarization will point in the same direction as that of the spontaneous polarization of the AlGaN and GaN.\nThe advantage of using nonpolar planes over c-plane nitrides is that the total polarization will be reduced. There may even be zero polarization for specific alloy compositions on specific planes. Such scenarios will be discussed in detail in future scientific papers. The important point is that the polarization will be reduced compared to that of c-plane nitride structures.\nAlthough high performance optoelectronic devices on nonpolar on-axis m-plane GaN have been demonstrated, it is difficult to obtain long wavelength emission from InGaN/GaN MQWs grown on m-plane GaN. This is probably due to the low In incorporation of the InGaN/GaN MQWs. The emission wavelength of devices grown on m-plane is typically 400 nm, while the emission wavelength of devices grown on c-plane is 450 nm at the same growth condition. Reducing the growth temperature increases the In incorporation; however, crystal quality would be degraded. This would be a significant problem for applications such as blue, green, yellow, and white LEDs.\nThe present invention describes a technique for the growth of nonpolar III-nitride light emitting devices in which the emission wavelength from the devices can be controlled by the barrier thickness of the MQWs in the devices. For example, the present invention has obtained blue and green emission without the effect of polarization."}
{"text": "The present invention relates to satellite navigation systems. More particularly, the present invention relates to satellite navigation receivers and systems adapted for use in automatic landing of aircraft.\nGlobal navigational satellite systems (GNSS) are known and include the global positioning system (GPS) and the Russian global orbiting navigational satellite system (GLONASS). GNSS-based navigational systems are used for navigation and positioning applications. In the GPS navigational system, GPS receivers receive satellite positioning signals from a set of up to 32 satellites deployed in 12-hour orbits about earth and dispersed in six orbital planes at an altitude of 10,900 nautical miles. Each GPS satellite continuously transmits two spread spectrum, L-band signals: an L1 signal having a frequency fL1 of 1575.42 MHz, and an L2 signal having a frequency fL2 of 1227.6 MHz. The L1 signal from each satellite is modulated by two pseudo-random codes, the coarse acquisition (C/A) code and the P-code. The P-code is normally encrypted, with the encrypted version of the P-code referred to as the Y-code. The L2 signal from each satellite is modulated by the Y-code. The C/A code is available for non-military uses, while the P-code (Y-code) is reserved for military uses.\nGPS navigational systems determine positions by timing how long it takes the coded radio GPS signal to reach the receiver from a particular satellite (e.g., the travel time). The receiver generates a set of codes identical to those codes (e.g., the Y-code or the C/A-code) transmitted by the satellites. To calculate the travel time, the receiver determines how far it has to shift its own codes to match the codes transmitted by the satellites. The determined travel times for each satellite are multiplied by the speed of light to determine the distances from the satellites to the receiver. By receiving GPS signals from four or more satellites, a receiver unit can accurately determine its position in three dimensions (e.g., longitude, latitude, and altitude). A conventional GPS receiver typically utilizes the fourth satellite to accommodate a timing offset between the clocks in the receiver and the clocks in the satellites. Additional satellite measurements can be used to improve the position solution.\nModern aircraft are equipped with guidance equipment to enable automatic landings in low-visibility conditions. The equipment, which generates deviations from a defined path based on signals received from terrestrial radio-navigation aids, must be designed to ensure that it is extremely improbable that it provides misleading guidance data (i.e., the equipment must have high integrity). Stringent integrity requirements can be met with monitoring of receiver performance. These monitoring functions must detect faults that could cause misleading information while providing a low false alarm rate (i.e., high continuity).\nThe typical accuracy of satellite navigation systems is comparable to terrestrial radio-navigation systems. However, current satellite navigation receivers do not meet the stringent integrity and continuity requirements necessary for incorporation into aircraft automatic landing systems.\nA satellite navigational receiver landing system having retrofit compatibility with integrated landing system (ILS) receivers, in accordance with the invention, includes radio frequency circuitry which converts satellite signals into intermediate signals. A first processing channel coupled to the radio frequency circuitry generates a first position, velocity and time (PVT) solution as a function of the intermediate signals, and provides a PVT and guidance data output. A second processing channel coupled to the radio frequency circuitry generates a second PVT and guidance solution as a function of the intermediate signals. Monitor circuitry provides monitor output signals as a function of the first and second PVT and guidance solutions. Shutdown circuitry provides the first PVT and guidance solution to the PVT and guidance data output under the control of a shutdown signal from the monitor circuitry. Monitoring functions ensure PVT and guidance accuracy is maintained within specified limits."}
{"text": "Air distillation plants generally include adsorbers so as to subject the air, compressed on entering the plant, to a drying and decarbonation treatment before it is cooled and then introduced into the distillation column proper. For this purpose, the air to be treated flows under a pressure typically of about 5 bar absolute through a bed of an adsorbent based in the adsorber. Once the adsorption capacity of the bed has been exhausted, the adsorbent is regenerated by means of a dry and decarbonated gas. This gas may be an off-gas recovered from the distillation column and injected into the adsorber in the reverse direction. In order to increase the regeneration effectiveness, this off-gas is heated to a high temperature of about 200.degree. C. This type of adsorber, in which the adsorbent is subject to cyclic temperature variations, is known by the name \"TSA (Temperature Swing Adsorption) adsorber\".\nHeating the gas as well as the adsorbent represents a major part of the energy consumed by the distillation plant. In order to minimize this consumption, it is therefore necessary to reduce the heat exchange between the adsorber and its environment.\nThermal insulation means consisting of an outer covering of the walls of the adsorber are known. This type of insulation proves to be very effective for the entire adsorber but does have certain drawbacks when the adsorber is subjected in a cyclic manner to temperature variations. This is because, since the walls of the adsorbers are generally made of metal, these constitute, depending on the phase of the cycle a heat store or heat sink tending to delay establishment of a steady state after a temperature change. In the regeneration phase, the walls absorb part of the heat provided by the heated gas and thus delay the heating of the adsorbent, in particular in regions close to the walls. In the adsorption phase, they restore this heat, thereby decreasing the adsorption effectiveness of the adsorbent. The walls therefore act out of step with the cycles and considerably decrease the effectiveness of the adsorber.\nThis phenomenon could be counteracted by increasing the size of the adsorbers and consequently by reducing, for a constant flow, the frequency of the cycles. However, recent developments in distillation plants have led to a reduction in the size of the adsorbers so as to reduce the amount of adsorbent and the cost of an adsorber. Because of the fact that the amount of adsorbent is reduced, the frequency of the cycles is considerably increased.\nIn order to improve the effectiveness of the known absorbers, adsorbers are known which are equipped with internal thermal insulation. Such internal thermal insulation must satisfy particular requirements:\nit must not allow bypassing of gas between the wall of the adsorber and the bed of adsorbent, so that the gas that is to be treated and the gas that has been treated are always separate; PA1 it must withstand the repeated compressions and decompressions imposed by the operating cycle and be able to attain pressure variations of about ten bar; and PA1 it must easily allow periodic hydrostatic testing, this being generally required by the authorities in the operating countries. PA1 the blocking means include at least one separating wall fitted between the said wall and the shell; PA1 the adsorber comprises means for decreasing natural convection in the gas-filled cavity; PA1 the means for decreasing natural convection in the gas-filled cavity comprise a sheet, in particular a metal sheet, placed between the wall and the shell in the gas-filled cavity in order to obtain a double gas-filled cavity; PA1 the sheet is centered with respect to the said wall and to the shell; PA1 the gaps between the sheet and the wall and between the sheet and the shell are between 10 and 40 mm, and preferably equal to 20 mm; PA1 the absorber comprises elements, especially point or linear elements, forming spacers between, on the one hand, the wall and the sheet and, on the other hand, between the sheet and the shell; PA1 the elements forming spacers between the wall and the sheet are offset with respect to the elements forming spacers between the sheet and the shell; PA1 the elements forming spacers are indented portions and raised portions preferably produced by goffering in the sheet, these indented portions and raised portions having, in particular, pointed ends; PA1 the indented portions and raised portions of the sheet extend parallel to the axis of the shell, over the entire length of that dimension of the sheet which is parallel to this axis, and have a cross-section of approximately triangular shape; PA1 the indented portions and raised portions alternate on each face of the sheet; PA1 in the case of an adsorber having a shell which is extended into a part forming a dome, the adsorber includes at least one and preferably two layers of a felt of fibres applied against the inside of the dome; PA1 for an adsorber intended, in particular, to operate under a varying pressure, the felt of fibres includes ceramic fibres and has a ratio between the fluid resistivity and Young's modulus of elasticity of less than 1 s/m.sup.2 ; PA1 the thickness of the or each layer of felt is between 20 and 30 mm, and preferably equal to about 25 mm; PA1 the layers of felt are attached to the inside of the dome by means of removable attachment devices; PA1 the removable attachment devices comprise rods, fastened to the dome and intended to transpierce the layers of felt, and washers for holding the felt against the dome, the washers engaging with the rods, in particular by a clip-in mechanism; and PA1 the adsorber is used in particular for the purification of air sent into an air distillation plant, the active material being an adsorbent.\nDocument US-A-3,925,041 teaches a TSA adsorber equipped with internal thermal insulation. This adsorber comprises a cylindrical shell oriented vertically, the ends of which are closed by ellipsoidal parts forming, respectively, a dome and a bottom. A bed of adsorbents is placed in the cylindrical shell\nIn the region of the bed of adsorbent, the internal thermal insulation is produced by sheets of a rigid insulant, these being applied end to end against the shell and compressing, against this shell, a layer of a fibrous insulating material. This thermal insulation is therefore placed between the shell and the bed of adsorbent. In order to prevent bypassing of the gases through this insulation during operation, metal sheets are applied against the shell, covering the sides of the sheets of insulant so as to deflect the flow of gas towards the bed of adsorbent.\nFurthermore, a fibrous insulant is applied by means of bolts against the internal parts of the dome and of the bottom of the adsorber.\nAlthough this internal thermal insulation considerably increases the effectiveness of an adsorber, it is not suitable for undergoing a hydrostatic test. This is because, after such a test, the thermal insulation is wet and ineffective. It must therefore undergo lengthy, difficult and consequently very expensive drying. Moreover, fitting the sheets of insulant and the metal sheets in the region of the bed, of adsorbent proves to be difficult and lengthy so that only specialists are able to place this insulation against the internal walls of the adsorber. This known insulation therefore considerably increases the cost of manufacturing of the adsorber."}
{"text": "Serotonin has been implicated in a number of diseases, disorders, and conditions that originate in the central nervous system, including diseases, disorders, and conditions related to, for example, sleeping, eating, perceiving pain, controlling body temperature, controlling blood pressure, depression, anxiety, addiction and schizophrenia. Serotonin also plays an important role in peripheral systems, such as the gastrointestinal system, where it has been found to mediate a variety of contractile, secretory, and electrophysiologic effects.\nBecause of the broad distribution of serotonin within the body, there is a need for drugs that affect serotonergic systems. In particular, agonists, partial agonists, and antagonists of serotonergic systems are of interest for the treatment of a wide range of disorders, including anxiety, depression, hypertension, migraine, obesity, compulsive disorders, schizophrenia, autism, neurodegenerative disorders (e.g., Alzheimer's disease, Parkinsonism, and Huntington's chorea), and chemotherapy-induced vomiting.\nThe major classes of serotonin receptors (5-HT1—7) contain one to seven separate receptors that have been formally classified. See Glennon, et al., Neuroscience and Behavioral Reviews, 1990, 14, 35; and D. Hoyer, et al. Phannacol. Rev. 1994, 46, 157-203.\nFor example, the 5-HT2 family of receptors contains 5-HT2a, 5-HT2b, and 5-HT2c subtypes, which have been grouped together on the basis of primary structure, secondary messenger system, and operational profile. All three 5-HT2 subtypes are G-protein coupled, activate phospholipase C as a principal transduction mechanism, and contain a seven-transmembrane domain structure. There are distinct differences in the distribution of the three 5-HT2 subtypes in a mammal. The 5-HT2b and 5-HT2a receptors are widely distributed in the peripheral nervous system, with 5-HT2a also found in the brain. The 5-HT2c receptor has been found only in the central nervous system, being highly expressed in many regions of the human brain. See G. Baxter, et al. Trends in Pharmacol. Sci. 1995, 16, 105-110.\nSubtype 5-HT2a has been associated with effects including vasoconstriction, platelet aggregation, and bronchoconstriction, as well as certain CNS effects, while subtype 5-HT2c has been associated with diseases that include depression, anxiety, obsessive compulsive disorder, addiction, panic disorders, phobias, psychiatric syndromes, and obesity. Very little is known about the pharmocologic role of the 5-HT2b receptor. See F. Jenck, et al., Exp. Opin. Invest. Drugs, 1998, 7, 1587-1599; M. Bos, et al., J. Med. Chem., 1997, 40, 2762-2769; J. R. Martin, et al., The Journal of Pharmacology and Experimental Therapeutics, 1998, 286, 913-924; S. M. Bromidge, et al., 1. Med. Chem., 1998, 41,1598-1612; G. A. Kennett, Drugs, 1998, 1, 4, 456-470; and A. Dekeyne, et al., Neurophannacology,1999, 38, 415-423.\nWO 93/13105 discloses thiophene derivatives; U.S. Pat. No. 4,414,225 discloses thiophene, furan and pyrrole derivatives and WO 96/12201 discloses furan derivatives."}
{"text": "The present invention relates to a mechanical detent for holding a valve spool in a predetermined position along its longitudinal axis.\nOne application of a detent is to hold a spool that is spring-centered, i.e., the spool returns automatically to a neutral position when the detent is disengaged. Another application of a detent is to hold a spool that is not spring-centered, i.e, the spool is manually moved to all positions including neutral. In addition to physically holding the spool in a certain position (\"stop\" position), a mechanical detent may also be used for a \"feel\" position. To move the spool past a feel position requires the operator to exert a higher force, indicating that a particular stage of the spool stroke has been reached or passed through. However, the spool is not physically held in this position.\nMaintaining a valve spool in a predetermined position by means of a mechanical detent has been attempted in many forms. Various forms of the detent include spring-loaded radial detents and C-shaped spring detents which surround the circumference of the valve spool. All of the above detents engage a groove or a step on the periphery of the valve spool.\nThe spring-loaded radial detents engage a groove on the spool stem, producing a radial force at a single point of contact. Some problems with the spring-loaded radial detents are that they require a number of parts, take up a large amount of space in the housing of the valve spool, and are costly both in part costs and assembly time. For example, in the case of a ball-type radial detent, the ball construction must display a radius approximately two or more times greater than the depth of the groove in order to provide a means for effectively moving the ball out of the groove. Thus, the deeper the groove is made, the larger the ball size must be. In trying to keep the ball size within an acceptable size range, while maintaining an acceptable groove depth, this required ratio between the ball and the groove results in a relatively shallow groove and a relatively large ball.\nDue to the built-in requirement of this relatively shallow groove, this construction creates the need for heavy loading of balls by spring(s) in order to maintain the ball in the groove. Additionally, where the ball size must be at least two times the groove depth, the only way to minimize the overall size of the valve structure is to have close toleranced housing valve spool and engaging parts. Such close toleranced parts increase the overall cost of the valve structure.\nIn the case of a C-shaped spring detent, the C-shaped spring surrounds the circumference of the valve spool and radially extends around the grooveless surface of the valve spool. When the groove of the valve spool is moved past the C-shaped spring, it springs radially inward, engaging the groove. The C-shaped spring provides essentially continuous contact around the surface of the spool. As in the case of the ball design, the C-shaped spring must have a relatively large cross-sectional diameter with respect to the groove depth since it must extend well beyond the surface of the spool to be effectively moved out of the groove. This required dimension of the groove engaging element creates the need for a relatively shallow groove to keep the C-ring size in an acceptable range and a corresponding heavy loading to maintain the element in the groove.\nIn addition to the constraint of requiring a relatively large groove engaging element, a relatively shallow groove, and a corresponding heavy load, the structure of the C-shaped spring and the ball-type detent also requires large spring size in order to have as low a rated spring (load/deflection) as possible. A low rated spring is necessary in order to have a maximum groove depth. Having a groove that is as deep as possible is desireable for manufacturing and tolerancing purposes. A low rated spring is necessary in the case of a maximum groove depth in order to move the element out of the groove. Without a low rated spring, a large load is required to move the element out of the groove. However, due to the construction of the above-mentioned detents it is necessary to further increase the overall size of the spring to achieve a low rated spring.\nFor example, to have a low rate C-shaped spring, the overall diameter D of the spring (the full length of the \"C\") must be large with respect to the diameter d of the wire of the spring, further increasing the size of the structure. The rate of the C-shaped spring is related to the ratio d/D. The lower this ratio, the lower the spring rate. Since the diameter d of the wire must be at least two times the groove depth, the numerator of this rate, d, has a somewhat \"fixed\" minimum value. Thus, to have a lower spring rate, the denominator, D, the overall size of the spring, must be large. Where a spool size is predetermined, modification of the spool may be necessary in order to accomodate a large, low rate spring.\nIn the case of the ball-type detent, having a coil spring having a wire size d and an overall spring length D, the spring rate is also related to the ratio d/D. To obtain a low spring rate, the length of the spring, D, must be large, making the overall size of the spring very large. The length of the spring especially affects the size of the total valve structure due to its orientation. In the ball-type detent, the coil spring is positioned radially outward from the ball. Thus, each increase in spring length further extends the radial distance between the spool and the housing required to accomodate the ball and spring. Although the surrounding parts can be constructed to be close-fitting to minimize the size, this is an expensive process.\nAnother disadvantage of the spring-loaded radial detents and the C-shaped spring detents is that, since the radius of the groove engaging element must be substantially greater than the depth of the groove, an enlarged load angle results. The load angle is defined as the angle between a horizontal axis through the groove engaging element and the axis traveling through the point of contact of the engaging element against the side wall of the groove and the center of the element. The load angle is equal to the angle of the side wall of the groove, or the groove angle. Due to the large load required in these types of detents, as well as the correspondingly high loading rates involved, these systems usually employ an edge contact point, whereby the outer diameter edge of the spool stem at the point where the groove begins contacts the wire or ball. This construction provides poor load angle control and high point contact stresses.\nExamples of the spring-loaded radial detents are found in U.S. Pat. Nos. 3,986,701 to Hopkins, 4,040,675 to Richmond et al. and 4,260,132 to Habiger. Examples of the C-shaped detent are shown in U.S. Pat. Nos. 4,413,805 to Green et al. and 4,185,661 to Gill et al.\nIt is an objective of the invention to create a detent having a groove engaging element having a radius which is independent of the groove size; the radius may be larger than, equal to or smaller than the depth of the groove.\nIt is another objective of this invention to create a detent having a groove engaging element which provides improved \"feel\" and \"stop\" characteristics with a reduced diameter.\nIt is another objective of the invention to create a detent which does not require a large number of parts and which requires reduced space in the valve spool assembly.\nAnother objective of this invention is to create a detent having improved load angle control and reduced contact stresses.\nAnother objective of this invention is to create a detent which provides improved \"feel\" and \"stop\" characteristics upon engagement and disengagement of the groove.\nAnother objective of the invention is to create a detent such that the clearance between the spool and the engaging elements is not as critical.\nAnother objective of the invention is to be able to place a detent anywhere on the spool and to avoid having to modify the spool to accomodate a low rate spring.\nIt is a further objective of the invention to create a detent which is easy to assemble and disassemble.\nThese objectives are achieved by a detent having a linear element spring having extensions which extend a sufficient distance beyond the surface of a spool to be engageable with the engaging element which moves the spring out of circumferential grooves on, the spool. Thus, the wire diameter of the spring is not critical, and may, for example, be reduced since the groove engaging portion of the wire is not the primary contact surface used to force the spring out of the groove.\nThe spring may be formed from a single wire having three or more linear segments interconnected at elbows. Alternatively, one or more linear springs may be positioned by a retainer. The extensions should extend a distance at least twice the diameter of the spring wire. Since the radius of the spring wire is not dependent on the particular desired groove depth, the radius of the spring wire can be greater than, equal to or less than the depth of the circumferential groove. Additionally, since there is no minimum required value for the spring wire in relation to the groove depth, a low rate spring may be designed without requiring a large spring size. Also, since the spring wire need not be at least two times the size of the groove depth, the point of contact between the groove engaging element and the side wall of the groove is below the edge point where the groove meets the outer diameter of the spool.\nFurther objects, features, and advantages of the present invention will become more apparent from the following description when considered in conjunction with the accompanying drawings which show, for purposes of illustration only, embodiments in accordance with the present invention."}
{"text": "1. Field of the Invention\nThe invention relates to a vehicle seat, especially a motor vehicle seat with a cushion support of the seat part, such that, the cushion support can be adjusted relative to the freely pivotable backrest by means of an adjustment device. As used herein, a freely pivotable backrest is understood to mean a backrest which can be tipped forward without having to activate the adjustment drive for adjusting the slant position.\n2. Description of the Prior Art\nWith the aforementioned vehicle seats, a relatively strong cushion pressure between the cushion of the backrest and the seat cushion is tolerated, at least when the latter is set into its topmost position by means of the height-and/or slant-adjustment device in the region of its rearward end. This situation is tolerated, despite its associated disadvantages, to prevent a considerable gap from existing between the rearward end of the seat cushion and the lower end of the backrest cushion when the seat cushion is in its lowermost position."}
{"text": "1. Field of the Invention\nThis invention relates to a process for repairing an electrical printed circuit and more particularly to repairing an electrical interruption in the circuit path.\n2. Information Disclosure Statement\nConductor path interruptions occur during manufacturing, repairing or as a result of physical shock or vibration. An interruption in one circuit connector will in many cases cause the total failure of the printed circuit board. Accordingly, a small defect in the printed circuit may necessitate the replacement of the entire printed circuit board and the associated components.\nOne process utilizes a device which holds a specifically shaped metallic part above the conductor path interruption. The device lowers the metallic part onto the conductor path and micro-resistance welds the shaped metallic part to the circuit so as to electrically reestablish the conductor path interruption.\nAnother prior art method seems to suggest the use of an insulated wire extending or bridging the electrical gap in the printed circuit.\nIn another process for manufacturing printed circuits, a resinous conductive ink circuit is printed on an electric board. While the ink is wet, the inked circuit is covered with a conductive metal powder. The powder is pressed into the ink and is allowed to cure. A solder paste is printed over the circuit and the board heated to cause the solder to alloy with the ink and powder to complete the method. A particular apparatus is employed to carry out the process.\nAnother process discloses an electrically conductive adhesive for bonding small semi-conductor components to a substrate. This patent fails to teach or suggest any tendency for a substance to bridge any gap. The disclosure of this patent is limited to the adhering of one small object to another through an electrically conductive bond.\nAnother method is drawn to an improved composition for a metallic ink as well as an improved method of use involving a low melting point metal first used in conjunction with a conductive ink. The combination is fired after the ink is screened in place in order to produce a relatively thick film on a substrate or semi-conductor device for making subsequent electrical contact by any one of various methods including clamping and soldering.\nAnother method utilizes a rather complex system to find a conductor path break and suggests the the gap may be bridged by the use of an insulated wire extending point to point across the gap.\nAnother method involves making contact between two conductors of printed circuits on opposite sides of a non-conducting substrate used as a circuit board. Holes are drilled through the board at the appropriate points. In the preferred embodiment, the sides of the holes are electroplated with a conductive material to permit electrical communication from circuits on one side of the board to the other.\nThe aforementioned methods utilize rather complex devices or insulated wire to effect the repair of a conductor break.\nAccordingly, it is an object of this invention to provide a method of repairing an electrical conductor path interruption on a printed circuit board which is economical and simple to perform.\nIt is a further object of this invention to provide a method of repairing an electrical conductor path interruption on a printed circuit board in the presence of heat sensitive electrical components.\nIt is a further object of this invention to provide a method of repairing an electrical conductor path interruption on a printed circuit board without the need for welding a metallic conductor at the site of the interruption.\nIt is a further object of this invention to provide a method of repairing an electrical conductor path interruption on a printed circuit board without the need for metallic strips to bridge the interruption.\nIt is a further object of this invention to provide a process which allows for a simple and complete repair of conductor path breaks without the use of a complex device.\nIt is a further object of this invention to provide a method of repairing an electrical conductor path interruption on a heat sensitive dielectric printed circuit board such as impregnated paperboard.\nThe foregoing has outlined some of the more pertinent objects of the present invention. These objects should be construed to be merely illustrative of some of the more pertinent features and applications of the invention. Many other beneficial results can be obtained by applying the disclosed invention is a different manner or modifying the invention within the scope of the disclosure. Accordingly, other objects and a fuller understanding of the invention may be had by referring to the summary of the invention and the detailed description describing the preferred embodiment in addition to the scope of the invention defined by the claims taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates to an apparatus generating and selecting ions used in a heavy ion cancer therapy facility according to independent claims.\nFrom U.S. Pat. No. 4,870,287 a proton beam therapy system is known for selectively generating and transporting proton beams from a single proton source. The disadvantage of such a system is, that the flexibility to treat patients is quite limited to relatively low effective proton beams.\nIt is an object of the present invention to provide an improved apparatus for generating and selecting different ions useful in an ion beam cancer therapy facility.\nThis object is achieved by the subject matter of independent claim 1. Features of preferred embodiments are defined by dependent claims.\nAccording to the invention an apparatus is provided for generating, extracting and selecting ions used in an ion cancer therapy facility. The apparatus comprises an independent first and an independent second electron cyclotron resonance ion source for generating heavy and light ions, respectively. Further is enclosed a spectrometer magnet for selecting heavy ion species of one isotopic configuration positioned downstream of each ion source; a magnetic quadrupole triplet lens positioned downstream of each spectrometer magnet; a switching magnet for switching between high-LET ion species and low-LET ion species of said two independent first and second ion sources. An analyzing slit is located at the image focus of each spectrometer magnet and a beam transformer is positioned in between the analyzing slit and the magnetic quadrupole triplet.\nSuch an apparatus has the advantage, that the possibility to help patients is largely improved by providing two independent ion sources and a switching magnet to select the proper ion species for an optimal treatment. Further the apparatus according to the present invention has the additional advantage that two independent spectrometer lines (one for each ion source) increase the selectivity of the apparatus and improve the purity of the ion species by separating with high accuracy the ion species selected for acceleration in the linac from all the other ion species extracted simultaneously from the ion sources.\nFor the intensity controlled rasterscannner ion beam application system different beam intensities within an intensity range of 1/1000 are provided in a preferred embodiment of the invention for each individual synchrotron cycle. The apparatus according to the present invention has the advantage to control the beam intensity at a low energy level in that the beam is destroyed along a low energy beam transport (LEBT) line in between the magnetic quadrupol triplet and an radio frequency quadrupole accelerator (RFQ). In particular, irises with fixed apertures are provided after a switching magnet as well as before and after a macropole chopper and at an RFQ entrance flange. An intensity measurement of the relative intensity reduction versus the magnet current of the center quadrupole of the magnet quadrupole triplet lens downstream of the image slit of the spectrometer is carried out for the apparatus of the present invention and shows that the beam intensity is reduced by about a factor of 430 starting from the default setting of the quadrupole magnet down to zero current. A further reduction of the beam intensity leading to a degradation factor of 1000 can be achieved by an additional reduction of the field of the third quadrupole of the magnetic quadrupole triplet. A very smooth curve is obtained, providing a good reproducibility of the different intensity levels.\nTherefore, the present invention avoids unnecessary radioactive contamination of the machine since beam intensity is controlled at the lowest possible beam energy, i.e. in said low energy beam transport line. Because the synchrotron injection scheme is not changed for the different beam intensity levels, i.e. the number of turns injected into the synchrotron are the same in all cases, the full dynamic range of 1000 is provided by the intensity control scheme in the LEBT according to the present invention. In the apparatus of the present invention the beam loss occurs mainly in the LEBT, i.e. the relative intensity reduction is almost the same measured directly behind the LEBT at a low energy level and measured in the Therapy beam line at an high energy level.\nFurthermore, beam profiles are measured at different locations along the accelerator chain and at the final beam delivery system of the therapy beam line. No differences could be observed in the beam profiles as well as in the beam positions for the different beam intensities. This is a very important advantage of the present invention in order to provide reliable and constant and not intensity dependent beam parameters at the treatment locations particularly when the when the apparatus of the present invention is applied for a heavy ion cancer therapy facility.\nThe beam transformer positioned in between the analyzing slit and the magnetic quadrupole triplet has the advantage to measure and monitor one-line the ion beam current of the ion species selected for acceleration without destroying the ion beam. Because this transformer is positioned upstream of the magnetic quadrupole triplet used for the intensity reduction the beam transformer monitors continuously the non-degraded ion beam current while intensity of the linear accelerator beam can be changed from pulse to pulse using triplet magnets. This is very important for an on-line monitoring of the performance of the selected ion source.\nIn a first preferred embodiment a solenoid magnet is located at the exit of each ion source. This embodiment of the present invention has the advantage that the ion beams extracted of each ion source are focused by a solenoid magnet into the object point of the spectrometer.\nIn an other preferred embodiment a magnetic quadrupole singlet is positioned downstream of each ion source. This quadrupole singlet has the advantage to increase the resolution power of each spectrometer system and to provide a flexible matching between the ion sources and the spectrometer systems.\nIn a further preferred embodiment the ion sources comprise exclusively permanent magnets. These permanent magnets provide a magnetic field for the ion sources and have the advantage that no magnet coils are required, which would have a large power consumption for each ion source. Additionally to the large power consumption these magnet coils have the disadvantage, that they need a high pressure water cooling cycle, which is avoided in the case of permanent magnets within the ion sources of the present invention. This has the advantage to reduce the operating costs and increase the reliability of the apparatus of the present invention.\nA further preferred embodiment of the present invention comprises beam diagnostic means which are located upstream each spectrometer magnet. Such beam diagnostic means can measure the cross-sectional profile of the beam and/or the totally extracted ion current. Said beam diagnostic means preferably comprises profile grids and/or Faradays cups.\nA further embodiment of the present invention provides a beam diagnostic means located at each image slit. This embodiment has the advantage to measure the beam size and beam intensity for different extracted ion species and to record a spectrum.\nIn a preferred embodiment of the invention, said focusing solenoid magnet is positioned downstream of said macropulse chopper and upstream of said radiofrequency quadrupole accelerator. This has the advantage that the beam is focused by the solenoid magnet directly to the entrance electrodes of the radio frequency quadrupole within a very short distance between the solenoid lens and the beginning of the RFQ electrodes of about 10 cm.\nA further preferred embodiment of the present invention provides diagnostic means comprising a Faraday cup and/or profile grids within the low energy beam transport system (LEBT) downstream of a switching magnet. These diagnostic means are not permanently within the range of the ion beam, but are positioned into the range of the ion beam for measurement purposes. The Faraday cup captures all ions passing the switching magnet and the profile grids measure the local distribution of ions within the beam cross section. During an operation cycle these diagnostic means are driven out of the range of the ion beam.\nIn a further preferred embodiment of the present invention the alternating stems within said radio frequency quadrupole are mounted on a common water cooled base plate. This has the advantage that the energy loss of the RFQ is conducted toward to outside of the chamber and do not damage the stems or the electrodes of the RFQ.\nIn a further preferred embodiment of the present invention the base plate is made of an electrical insulating material. This has the advantage that the stems are not short circuit, though they are acting as inductivity whilst said mini-vane pairs forming electrodes are acting as capacitance for a xcex/2 resonance/structure."}
{"text": "1. The Field of the Invention\nThe invention relates to the structure of and methods of presetting bellows-type gas lift valves.\n2. The Prior Art\nA wide variety of bellows-type gas lift valves is known in the art, many of which include a fixed mechanical stop to limit the valve stem travel. Such a stop is desirable to prevent damage to the bellows in the event that an excessively high operating pressure is applied to the valve.\nOne example of a bellows-type gas lift valve having a mechanical stop is illustrated in U.S. Pat. No. 3,192,869, issued July 6, 1965, to McCarvell et al.\nOther gas lift valves, such as the Camco type \"BK\" bellows-type valves, also have a mechanical stop for limiting the valve stem travel. In this type of valve, the stop is positioned prior to welding the bellows into the valve, and no provision is made for repositioning the stop after the bellows is affixed, i.e., during the bellows presetting operation.\nIt is known in the art to preset the bellows of a gas lift valve after assembly by overpressuring the valve. For example, a valve which is to have an operating pressure of 600-800 psi may be overpressured to, for example, 3500 psi. The purpose of presetting the bellows is to permanently deform the bellows so that the valve will open at the desired operating pressure. The purpose of presetting the bellows is to assure the operating pressure will not change if an excessive pressure is encountered in the well, say as high as 3500 psi.\nSince it appears that the conventional presetting procedure described above does not cause sufficient compression of the bellows, the bellows will \"stack up\" at operating pressure, thereby limiting the stem travel before the mechanical stop is reached. Moreover, it has been found that stem travel is non-linear with applied pressure, and it would be desirable to make the relationship more nearly linear. Still further, the prior art valve assemblies and presetting methods do not provide entirely consistent operating pressure/stem travel characteristics from valve to valve. That is, some valves will fully open at a pressure differential which differs significantly from that required by other valves of the same construction and preset by the same process."}
{"text": "The present invention relates to a vehicle windshield wiper assembly.\nReferring to FIGS. 6 and 7, a conventional windshield wiper assembly 200 typically includes a wiper motor 210, one or more wiper blades secured to wiper arms, and a pivotable linkage mechanism therebetween.\nThe wiper motor 210 has a housing 211 and a motor 220 which is connected to the housing 211. The housing 211 has a top surface 212 and a side surface 213. A projection 214 extends from the top surface 212. The projection 214 pivotally supports an output shaft 230 of the wiper motor 210. When the wiper motor 210 is driven, the output shaft 230 rotates in a predetermined direction r1.\nThe wiper motor 210 has a mounting bracket 215 integrally formed with the side surface 213. The mounting bracket 215 is attached to a bar when mounted within a vehicle. The wiper assembly 200 is mounted to the vehicle at mounting parts 250. A link mechanism 260 which is attached to the output shaft 230 moves back and forth in predetermined ranges r2, r3 and r4 due to rotation of the wiper motor 210. As a result, the wiper arms reciprocate along a windshield.\nWhen the wiper motor 210 rotates while under the influence of obstacles such as snow, sleet, and the like at or near the end-of-stroke reverse positions, the load on the link mechanism 260 becomes high. As a result, each part of the wiper assembly 200 undergoes considerable stress. In particular, the housing 211 of the wiper motor 210 becomes highly stressed at the top surface 211, the side surface 212 and the mounting bracket 215."}
{"text": "Terminology:\nIn the following, references to skates, blades, or skate blades apply also to bobsled and iceboat runners unless otherwise noted. Sharpener or sharpeners will refer to small hand-held devices used to maintain or create sharp-edged ice-contact surfaces of skates in a hollow (concave) or flat shape by sliding the device back and forth along the blade. Abrader refers to the abrasive or cutting item (file, emery-cloth, stone, etc.) that grinds or cuts the blades. Convex and concave, which refer to the shape of the abrader or blade hollow, are defined as \"radius\" or \"radii\" in decimal-inch values.\nVariables that determine the radius values needed:\n1. In hockey, the forward requires great maneuverability for very sharp starts, turns, and stops; the goalkeeper requires great stability and side-to-side stepping/sliding action. Smaller radius for the former player, larger radius or flat for the latter. PA0 2. Speed skaters need a flat or large radius for maximum speed without sharp turns and with small variations for indoor smaller-turn rinks. PA0 3. Figure skaters may have several pairs of skates with a variety of radii for success in \"compulsory figures\" and the strenuous jumps and spins of the axels and lutzs. PA0 4. Added to the above are skaters' individual preferences as well as their choices depending on ice conditions (hard/soft, smooth/rough). PA0 A. Must have a self-aligning or simple-adjustment means to locate the center of the skate blade exactly in the center of the convex abrader. PA0 B. If the abrader is a long round file, stone, or cylinder, etc., that must be rotated after a sharpening to obtain a fresh, true convex radius, the worn surface must not affect the exact centering as stipulated in Specification A. PA0 C. Must be able to produce sharpened blade hollows of any radius value from 0.032\" to the infinite value of flat. PA0 D. Abrader components or surfaces that become worn out in sharpening should be available at low cost from local hardware or sports stores to compete easily with high-cost, inconvenient commercial machine grinding. PA0 E. Guides for the skate blades must be easily and accurately adjusted to accommodate all skate blade widths and to allow for wear of the guides. PA0 F. Simple, small pocket tools, such as nail files, folding knives, and screwdrivers, should be the only equipment needed to adjust, maintain, and replace parts. PA0 G. The sharpener should be light in weight. PA0 H. Sharpeners should not have sharp, protruding parts to damage clothing or harm the user. PA0 I. Parts that may come loose or lost--like screws, nuts, or washers--should be standard hardware items. PA0 J. The number of accurate sharpenings possible before replacement of the abrader or blade guides should be high for convenience and cost efficiency.\nOptimum Sharpener Specifications\nBased on the above many variables, the following describe the optimum requirements for all skate sharpeners:"}
{"text": "1. Field of the Invention\nThis invention relates generally to apparatus for use with the welding of piping to be used in nuclear power plants and systems. It relates specifically to the apparatus used for heat sinking and purging of the piping.\n2. Description of the Prior Art\nA common problem with known devices used in the welding of pipes is that they do not offer the flexibility of both cooling as a heat sink together with purging features. Known type devices merely provide heat sink or cooling functions, or purging functions, but not the combination of same.\nAnother problem with known type devices are that they are not readily adaptable for large scale operations or mass welding of numerous piping. That is, they are so complex and involved that much time is spent in setting them up and they are not easily transferable from one set of pipes to be welded to another set.\nAnother problem with known type devices is that they are not readily adaptable for use with either gas type purging functions and/or liquid-type purging and cooling functions.\nKnown prior art patents which may be pertinent to this invention are as follows:\nU.s. pat. No. 2,470,744 -- May 17, 1949 PA1 U.s. pat. No. 2,819,517 -- Jan. 14, 1958 PA1 U.s. pat. No. 3,194,466 -- July 13, 1965 PA1 U.s. pat. No. 3,431,945 -- Match 11, 1969 PA1 U.s. pat. No. 3,779,068 -- Dec. 18, 1973 PA1 U.s. pat. No. 3,902,528 -- Sept. 2, 1975\nNone of these known prior art devices offers the new and unique features of the invention disclosed herein."}
{"text": "In a known connectors of this kind, one part of the housing has two U-shaped clamps, which can be swung with recesses provided in the legs onto pins of the other part of the housing and can be locked by means of expansions, formed in the rim surfaces of the recesses as detents, at these expansions. Apart from the fact the clamps can be operated only separately in succession, whereby it is possible to tilt the parts of the housing with respect to each other, the expansions acting as detents require a significant amount of energy to swing on or release the clamps, having a negative mechanical effect on pins and clamps. In addition, the detents provided at the clamps do not function reliably. In addition, this known connector is difficult to manipulate due to the independently uncontrollable swivelling of the clamps before or during the plugging operation of the parts of the housing."}
{"text": "When a semiconductor device is fabricated, a semiconductor die is typically wire bonded to a metal lead frame. In a conventional wire bonding assembly, a plurality of connected lead frames, each having a semiconductor die attached thereto, is inserted along a track. Each lead frame is passed under a hold down clamp which maintains the lead frame within a window which is cut through the clamp. After the lead frame is clamped, the semiconductor die is wire bonded to the lead frame. Typically, gold wire is used and the lead frame is a metal such as a nickel-iron alloy or copper. Due to a variety of factors, such as the presence of impurities on bonding pads, some wires do not stick to the bonding pads. In order to further enhance the strength of the wire bond, the semiconductor die is typically heated from below the lead frame. A disadvantage with heating the semiconductor die from beneath the lead frame is that the die is generally too thick for enough heat to reach the bonding surface to improve the strength of the wire bond. When inadequate heat occurs, more wire bonding failures occur. In conventional die bonding assemblies, others have used hot gas chambers. However, gas chamber heating is inefficient and difficult to control when used with high speed bonding and ceramic capillary gas tubes."}
{"text": "This invention relates to serially connected logic units wherein each logic unit, such as an AND gate or an OR gate, is composed of a resistor connected serially to a combination of plural transistors connected in parallel and, more particularly, to the construction of one of the logic units of PNP transistors and another of the logic units of NPN transistors wherein the voltage drop across a transistor of a first of the logic units establishes an operating voltage level for feeding current to the transistor of a following logic unit.\nThere are increasing demands for higher performance of computers, data processing equipment, and control systems which employ logic circuits, the. higher performance being in terms of greater speed of operation Due to the wide use of logic circuits in the foregoing equipments, an increase in the speed of response of the logic circuit provides a significant improvement in the performance of the equipment.\nA typical form of construction of a logic circuit, such as an AND gate or an OR gate has an interconnection of a plurality of identical transistors in a parallel circuit combination, the combination then being connected serially with a resistor. The ]unction of the terminals of the parallel combination with a terminal of the resistor serves an output terminal of the gate, the opposite terminal of the resistor being connected to one terminal of a power supply and the opposite terminal of the parallel combination being connected to the other terminal of the power supply. With such a circuit arrangement, the transistors force current through the resistor in one direction for a fast response to the waveform of a pulse signal, a change of current in the opposite direction being accomplished relatively slowly in the manner of a capacitive discharge via the resistor with the capacitance being provided by stray capacitance of the circuit With such construction, the speed of response can be increased by reducing the resistance of the resistor.\nHowever, a problem arises in that the reduction of the resistance increases the dissipation of power within the logic unit. However, computers, data processing equipment, and control systems are limited in the amount of power which can be dissipated therein. This limitation, in turn, limits the amount of reduction in resistance which is permissible before exceeding power dissipation limits. As a result, there is a need for still further modification of the circuitry of a logic unit to produce the advantage of increased speed."}
{"text": "Implantable electrical stimulation systems have proven therapeutic in a variety of diseases and disorders. For example, spinal cord stimulation systems have been used as a therapeutic modality for the treatment of chronic pain syndromes. Peripheral nerve stimulation has been used to treat chronic pain syndrome and incontinence, with a number of other applications under investigation. Functional electrical stimulation systems have been applied to restore some functionality to paralyzed extremities in spinal cord injury patients.\nStimulators have been developed to provide therapy for a variety of treatments. A stimulator can include a control module (with a pulse generator), one or more leads, and an array of stimulator electrodes on each lead. The stimulator electrodes are in contact with or near the nerves, muscles, or other tissue to be stimulated. The pulse generator in the control module generates electrical pulses that are delivered by the electrodes to body tissue.\nConventional implanted electrical stimulation systems are often incompatible with magnetic resonance imaging (“MRI”) due to the large radio frequency (“RF”) pulses used during MRI. The RF pulses can generate transient signals in the conductors and electrodes of an implanted lead. These signals can have deleterious effects including, for example, unwanted heating of the tissue causing tissue damage, induced currents in the lead, or premature failure of electronic components."}
{"text": "The invention relates generally to an exhaust fan and more particularly to an axial-flow exhaust fan for use as a bilge blower in marine environments.\nTwo industrial forms of fans or blowers predominate, axial-flow fans and centrifugal, or radial-flow, fans. In centrifugal fans, air flows through the fan wheel in a mostly radially outward direction, while air flows in an axial-flow fan in an axial direction with almost no radial component.\nAxial-flow fans operate by deflecting axially directed air on airfoils, or blades. This deflection causes the air flow to take on a helical flow pattern past the airfoils. This flow shape has two flow components, tangential velocity and axial velocity. Of the two flow components, the axial velocity is the more important component for moving air through the fan. Guide vanes positioned either upstream or downstream of the airfoils serve to translate the tangential velocity component of the air flow into the axial velocity component.\nThere are two methodologies for determining the size, dimensions, and number and positioning of blades for an axial-flow fan. One method is testing a first axial-flow fan design to ascertain the air volume and static pressure of the fan. Rarely does a first design meet the desired running parameters, and thus redesigning one or more times becomes necessary. Redesigning costs man-hours, and often the result is a fan which is larger than originally anticipated that runs at higher speeds and consumes more brake horsepower than needed. In addition, redesigning often leads to uneven and turbulent air flow and to the creation of stalling effects in certain parts of the blades.\nA second methodology, which improves over the first methodology, is to design the axial-flow fan based upon desired outcome parameters as well as desired structural and design parameters. Such parameters may include high efficiency and low sound output over a wide range of operation, non-overloading brake horsepower, a steep pressure curve (little variation in air delivery), a large free delivery of air, large pressure safety margin, and compactness.\nOne problem experienced with conventional axial fans used in marine environments is that the fan housings are subjected to physical and thermal stresses which may alter the diameter of the housings.\nThe invention provides an axial-flow fan which includes a plurality of airfoils extending from a rotatable fan wheel hub, a motor engaged with the fan wheel hub through a shaft, and a housing into which the fan wheel hub and the motor are positioned. The housing has a circumference with a diameter and at least one supporting element surrounding the circumference of the housing. The supporting element inhibits variation in the diameter of the circumference of the housing.\nThe invention further provides a method of manufacturing an axial-flow fan. The method includes the steps of engaging a fan wheel hub with a motor through a shaft, mounting at least one straightening vane from an inner surface of a housing, positioning the fan wheel hub and the motor within the housing such that the motor is mounted on one of the straightening vanes, and surrounding the housing with at least one supporting element which inhibits variation in the diameter of the circumference of the housing.\nWith these and other objects, advantages and features of the invention that may become hereinafter apparent, the nature of the invention may be more clearly understood by reference to the following detailed description of the invention, the appended claims and to the several drawings attached herein."}
{"text": "The present invention relates to the formation of a borophosphosilicate glass (xe2x80x9cBPSGxe2x80x9d) layer during the fabrication of integrated circuits on semiconductor wafers. More particularly, the present invention relates to an improved reflow process that reduces the thermal budget of a fabrication process while providing gap-filling properties that enable the BPSG layer to meet the requirements of modern day manufacturing processes.\nSilicon oxide is widely used as an insulating layer in the manufacture of semiconductor devices. A silicon oxide film can be deposited by thermal chemical vapor deposition (CVD) or plasma enhanced chemical vapor deposition (PECVD) processes from a reaction of silane (SiH4), tetraethoxysilane (Si(OC2H5)4), hereinafter referred to as xe2x80x9cTEOS,xe2x80x9d or a similar silicon-containing source, with an oxygen-containing source such as O2, ozone (O3), or the like.\nOne particular use for a silicon oxide film is as a separation layer between the polysilicon gate/interconnect layer and the first metal layer of MOS transistors. Such separation layers are referred to as premetal dielectric (PMD) layers because they are typically deposited before any of the metal layers in a multilevel metal structure. In addition to having a low dielectric constant, low stress and good adhesion properties, it is important for PMD layers to have good planarization and gap-fill characteristics.\nWhen used as a PMD layer, the silicon oxide film is deposited over a silicon substrate having a lower level polysilicon gate/interconnect layer. The surface of the silicon substrate may include isolation structures, such as trenches, and raised or stepped surfaces, such as polysilicon gates and interconnects. The initially deposited film generally conforms to the topography of the substrate surface and is typically planarized or flattened before an overlying metal layer is deposited.\nOne method developed to fill the gaps and xe2x80x9cplanarizexe2x80x9d or xe2x80x9cflattenxe2x80x9d the substrate surface involves forming a layer of relatively low-melting-point silicon oxide and then heating the substrate sufficiently to cause the layer to melt and flow as a liquid, resulting in a flat surface upon cooling. Such heating can be performed using either a rapid thermal pulse (RTP) method or conventional furnace, for example, and can be done in a dry (e.g., N2 or O2) or wet (e.g., steam H2/O2) ambient. Each process has attributes that make that process desirable for a specific application.\nBecause of its low dielectric constant, low stress, good adhesion and gap-fill properties and relatively low reflow temperature, borophosphosilicate glass (xe2x80x9cBPSGxe2x80x9d) is one silicon oxide film that has found particular applicability in applications that employ a reflow step to planarize PMD layers. Standard BPSG films are formed by introducing a phosphorus-containing source and a boron-containing source into a processing chamber along with the silicon- and oxygen-containing sources normally required to form a silicon oxide layer.\nAs semiconductor design has advanced, the feature size of the semiconductor devices has dramatically decreased. Many integrated circuits (ICs) now have features, such as traces or trenches that are significantly less than a micron across. While the reduction in feature size has allowed higher device density, more chips per wafer, more complex circuits, lower operating power consumption, and lower cost, the smaller geometries have also given rise to new problems, or have resurrected problems that were once solved for larger geometries.\nOne manufacturing challenge presented by submicron devices is minimizing the overall thermal budget of the IC fabrication process in order to maintain shallow junctions and prevent the degradation of self-aligned titanium silicide contact structures, among other reasons. Hence, for at least this reason, it is desirable to provide methods of forming planarized insulating layers, such as BPSG layers, with lower thermal budget requirements.\nThe present invention provides exemplary methods, apparatus and systems for planarizing an insulating layer, such as a borophosphosilicate glass (BPSG) layer or an undoped silicate glass (USG) layer, deposited over a substrate.\nIn one embodiment, the method includes loading a substrate having a BPSG layer deposited thereover into a substrate processing chamber. In one embodiment, the BPSG layer is a premetal dielectric (PMD) layer, although the BPSG layer may be positioned elsewhere in the circuit device within the scope of the present invention. The BPSG layer has an upper surface that is generally non-planar. The substrate is exposed to an ultraviolet (UV) light at conditions sufficient to cause a reflow of the BPSG layer so that the upper surface is generally planar. The UV light is produced with a UV lamp, a laser, other provided UV light sources, and the like. In this manner, photonic energy is used instead of thermal energy to cause the insulating layer to reflow. The reflow fills the gaps, vias, trenches and the like, producing a generally planar surface.\nIn one embodiment, the UV light has a wavelength of about 150 nmxc2x150 nm, although wavelengths throughout the UV spectrum may be used within the scope of the present invention. In alternative embodiments, the UV light has an energy level that is greater than about 10 electron volts (eV), and is about 15 eV.\nThe substrate, in one embodiment, is exposed to UV light for between about thirty (30) seconds and about fifteen (15) minutes. In another embodiment, the exposing step is maintained using UV light at an energy level that is at least about 10 eV, and for a duration that is at least about 30 seconds to produce sufficient reflow of the BPSG layer. In another embodiment, the UV light has a wavelength that is at least about 150 nm and the exposing step duration is at least about 30 seconds. It will be appreciated by those skilled in the art that the exposure time will depend, in part, on the UV light wavelength and/or energy, and the type and/or thickness of the insulating layer, among other things.\nIn one embodiment, the method includes maintaining a temperature in the substrate processing chamber between about 20 degrees Celsius and about 100 degrees Celsius during the exposing step. In this manner, a low thermal budget is used for the insulating layer reflow process.\nIn another embodiment, the exposing step exposes the substrate to UV light having a desired wavelength and a desired energy level to break at least some SiOH bonds in the BPSG layer. In this manner, the hydrogen content in the BPSG is reduced. Similarly, exposing the substrate to UV light helps densify the BPSG layer.\nIn another embodiment of the present invention, a method of forming a planarized insulating layer includes providing a substrate having a non-planar upper surface and depositing an insulating layer over the upper surface. The insulating layer has a generally non-planar upper surface, typically similar in contour to the substrate upper surface. The method includes exposing the insulating layer to a UV light at conditions sufficient to cause the insulating layer to reflow so that the insulating layer upper surface is generally planar. In one embodiment, the insulating layer comprises borophosphosilicate glass (BPSG), although other insulating layers, including other silicon oxide layers may be used within the scope of the present invention.\nIn one embodiment, the insulating layer is deposited by inserting the substrate into a substrate processing chamber and introducing a phosphorus-containing source and a boron-containing source into the processing chamber to deposit the BPSG insulating layer over the substrate. Examples of phosphorus-containing sources for use with the present invention include triethylphosphate (TEPO), triethylphosphite (TEPi), trimethylphosphate (TMOP), trimethylphosphite (TMPi), and similar compounds. Examples of boron-containing sources for use with the present invention include triethylborate (TEB), trimethylborate (TMB), and similar compounds.\nIn one embodiment, the depositing and exposing steps are both performed in a substrate processing chamber. Alternatively, the depositing and exposing steps are performed in separate processing chambers.\nIn one embodiment of the invention, a substrate processing apparatus comprises a processing chamber and a substrate holder located within the chamber for holding a substrate. A UV light source is coupled to the processing chamber and disposed to transmit a UV light towards the substrate holder, and hence towards the substrate. A controller for controlling the UV light source is included, with a memory coupled thereto. The memory includes a computer readable medium having a computer readable program embodied therein for directing operation of the UV light source. In one embodiment, the computer readable program includes a first set of instructions for controlling a wavelength of UV light produced by the UV light source, and a second set of instructions for controlling a duration that the UV light source produces UV light.\nIn another embodiment, the computer readable program further includes a third set of instructions for controlling an energy level of the UV light produced by the UV light source. In some embodiments, the apparatus or system further includes a gas distribution system coupled to the processing chamber for the deposition of an insulating layer on the substrate. In this manner, the same chamber can be used for both insulating layer formation and reflow.\nThese and other embodiments of the present invention along with many of its advantages and features are described in more detail in conjunction with the text below and attached figures."}
{"text": "1. Field of the Invention\nThe present invention relates, in a processing system having a plurality of CPUs each executing an online control program, to a method of recovering the failure of the online control program.\n2. Description of the Prior Art\nA prior-art method of recovering the failure of an online control program at high speed is called the \"hot standby recovery method.\" With this method, in correspondence with the online control program, a standby program is previously prepared in a state in which the process of the preceding stage of a recovery process for the failure has been completed, and the recovery process is performed with the standby program at the failure of the online control program, thereby to shorten the recovery period of time.\nThe official gazette of Japanese Patent Application Laid-open No. 11457/1984, etc. are relevant to the technique of this type.\nThe hot standby recovery method in the prior art has had the following problems since the redundant construction in which the standby program is prepared separately from the online control program is employed, at the failure of the online control program, a service cannot be started until the recovery process with the standby program is completed. Further the online control program cannot be recovered at high speed in a case where it has failed after the failure of the standby program."}
{"text": "The present invention relates generally to the field of data reduction, and more particularly to data reduction using a data analysis methodology.\nThe fabrication of printed circuit (PC) boards and related electrical equipment requires the diagnostic testing of performance factors or parameters including voltage and signal deviations, jitters, voltage rise and fall time and other signal integrity factors. The testing of these parameters helps identify anomalous data points from which a problem or failure in the PC board can be determined. Known testing methods of signal integrity in the manufacture of PC boards includes identifying probe or node points from a schematic or CAD layout of the PC board. These node points are located throughout the entire PC board in areas surrounding components, drivers, voltage points and the like and represent points of contact on the PC board.\nThe number of node points incorporated throughout a PC board can be large; consequently, raising the level of complexity required to perform any type of statistical or data related analysis of the testing results (raw data) obtained from the PC board. Additionally, these results include varied unit designations, such as, voltage, amperage, inductance, capacitance and time readings which further complicate data analysis of the results obtained from the diagnostic testing of the node points.\nOnce the node points of the PC board are tested and the results (raw data) are obtained, statistical analysis of the data is required in order to ascertain any problem or anomalous data which may indicate problem or defective areas in the PC board. Further, the large amounts of raw data often must be statistically processed in order to emphasize a small subset of data that may be useful to engineering personnel using the data to assess behavioral characteristics of circuitry on the PC board. To this end, there are many methods of performing statistical analysis on data including multiple regression and factor analysis. In using multiple regression, see, for example, Engineering Statistics by Douglas C. Montgomery et al., no ideal data model can be extrapolated due to the varied units which prevents quantitative aggregation of the raw data in a useful model form. The use of other known statistical factor analysis methods presents similar problems in consistently presenting guidelines for discerning problems at the node points.\nThe present invention provides a method and apparatus for performing data reduction to results of diagnostic testing of signal integrity of PC boards using a Rank-of-Ranks data analysis methodology. Raw data obtained from the diagnostic testing of signal integrity at node points situated on PC boards is manipulated to a succinct form using the Rank-of-Ranks data analysis methodology.\nAccording to the invention, the Rank-of-Ranks methodology includes an initial reduction of the quantity of analyzed node points on a PC board, by application of a node point reduction algorithm or selection criterion. The remaining node points to be tested (xe2x80x9cthe test nodesxe2x80x9d), each are identified by a node name and associated node position in (x-y-z) coordinate system. The test nodes are physically tested via a testing probe positioned on the end of a robot which moves along the node position coordinates. The testing at each node point includes sample measurements of selected testing parameters or factors. Standard Deviations and Mean are calculated for each selected testing parameter. All test data is compiled at a server and input into a spreadsheet for data manipulation using the Rank-of-Ranks methodology.\nData manipulation using the Rank-of-Ranks methodology according to the present invention includes calculating a Coefficient of Variation as a ratio of the Standard Deviation of a given factor, to the Mean, for each node point (and for each parameter). Each node point is sorted and ranked based on the Coefficient of Variation for each parameter at a time. At each node point, Rank values of all the parameters are summed and then sorted on the basis of the Rank values. Finally, the node points are Re-ranked based on the sorted Rank values. The Re-ranked node points are entered into an ordered list or matrix where critical node points are chosen from the top of the Re-ranked node points for study, re-design and/or further testing of the PC board.\nAdvantages of the present invention include reduction of node points involved in the testing of PC boards which allow the testing procedures to be more focused and to require less time thereby greatly increasing productivity and lowering manufacturing costs. Also, diagnostic testing of the PC boards can be targeted to include certain critical and control state node points."}
{"text": "Electrochemical biosensors are well known and have been used to determine the concentration of various analytes from biological samples, particularly from blood. Examples of such electrochemical biosensors are described in U.S. Pat. Nos. 5,413,690; 5,762,770; 5,798,031; and 6,129,823 each of which is hereby incorporated by reference.\nIt is desirable for electrochemical biosensors to be able to analyze analytes using as small a sample as possible, and it is therefore necessary to minimize the size of their parts, including the electrodes, as much as possible. As discussed below, screen-printing, laser scribing, and photolithography techniques have been used to form miniaturized electrodes.\nElectrodes formed by screen-printing techniques are formed from compositions that are both electrically conductive and screen-printable. Furthermore, screen printing is a wet chemical technique that generally allows reliable formation of structures and patterns having a gap width or feature size of approximately 75 μm or greater. Such techniques are well known to those of ordinary skill in the art.\nLaser scribing is a technique that usually uses a high power excimer laser, such as a krypton-fluoride excimer laser with an illumination wavelength of 248 nm, to etch or scribe individual lines in the conductive surface material and to provide insulating gaps between residual conductive material which forms electrodes and other desired components. This scribing is accomplished by moving the laser beam across the surface to be ablated. The scribing beam generally has a relatively small, focused size and shape, which is smaller than the features desired for the product, and the formation of the product therefore requires rastering techniques. Such a technique can be rather time consuming if a complex electrode pattern is to be formed on the surface. Further, the precision of the resulting edge is rather limited. This scribing technique has been used to ablate metals, polymers, and biological material. Such systems are well known to those of ordinary skill in the art, and are described in U.S. Pat. Nos. 5,287,451, 6,004,441, 6,258,229, 6,309,526, WO 00/73785, WO 00/73788, WO 01/36953, WO 01/75438, and EP 1 152 239 each of which is hereby incorporated by reference. It would be desirable to have a new method of forming electrodes which allows precise electrode edges, a variety of feature sizes, and which can be formed in a high speed/throughput fashion without the use of rastering."}
{"text": "Strong base anion resins are often used for decolorization purposes, in which colored impurities are removed from a feed stream. As one example from the cane sugar refining industry, decolorization is used to produce a crystallized sugar having an acceptable white color. Because many of the color bodies in a cane processing stream are hydrophobic and anionic, strong base anion resins are conventionally used for decolorization since the color bodies will bind to this type of resin through ionic and hydrophobic interactions. For this reason, cane refinery decolorization using strong base anion resins has become common.\nUsing ion exchange resin for decolorization has been less common in the beet sugar industry since the process, when operated conventionally, is capable of producing a white crystallized sugar without additional purification. However, more recent processes which increase the recovery of beet sugar, such as chromatography, have resulted in higher color feed to the crystallization step and have motivated a consideration of decolorization. The use of chemicals and waste disposal problems associated with anion resin decolorization have hindered widespread adoption of this technology. Similar problems exist in the cane sugar industry but are often accepted because decolorization is required to produce a sellable sugar.\nThe various problems associated with the anion resin decolorization process are primarily due to the method of regenerating the strong base anion resin. The regeneration is a well-known process and typically includes feeding a regenerant solution that includes sodium chloride (NaCl) and sodium hydroxide (NaOH) across the strong base anion resin, as shown in FIG. 1. The regenerant solution includes 10% NaCl and 1.0% NaOH. For additional regeneration, fresh NaCl and NaOH are added to the regenerant solution. The regenerant solution displaces the exchanged and adsorbed color bodies. The used regenerant is referred to as regenerant waste or spent regenerant, and is considered a waste stream. The spent regenerant contains most of the NaCl from the regenerant solution in addition to various organic color compounds previously removed by the strong base anion resin. Since the spent regenerant is a caustic brine, the spent regenerant is an undesirable material because it has a very high chemical oxygen demand (COD), high salt content, and poorly degradable color bodies. In some cases, the spent regenerant disallows the use of the process due to effluent discharge limits.\nTo overcome some of the waste problems, a reprocessing-recycle technique using nanofiltration to recover a fraction of the NaCl has been in use since the 1990s. In this method, the spent regenerant is passed through a crossflow membrane nanofiltration system, producing a membrane permeate and a membrane retentate, as shown in FIG. 2. A significant amount of the NaCl and NaOH of the spent regenerant passes through the membrane and is recovered in the membrane permeate. Fresh NaCl and NaOH are added to the membrane permeate for subsequent resin regeneration. A significant reduction in salt consumption may be achieved by recycling part of the spent regenerant in this manner. However, the process still requires a fresh NaCl and NaOH makeup and a spent regenerant bleed-off and salt/organic containing concentrate must be dealt with. The membrane retentate is sent to waste."}
{"text": "The present invention relates to wheelchairs, and particularly to wheelchairs capable of moving the seat and back for tilting and reclining. More particularly, the invention pertains to wheelchairs having legrests capable of being raised to support the legs of the wheelchair user.\nWheelchairs often have a fixed seat consisting of a seating surface and a back frame. The seating surface is usually either horizontal or slightly tilted back, with the front edge of the seating surface slightly higher than the rear edge of that surface. If the wheelchair user sits in the same position in a wheelchair for a long period of time, pressure is continuously applied to the tissue on the portion of the user\"\"s body (buttocks, legs, and/or back) that is bearing the user\"\"s weight in that position. Blood circulation to that tissue will be reduced, and ulcers or other problems can result.\nTo avoid these problems, it is necessary for people sitting in wheelchairs to shift their body weight from time to time. This is often accomplished by tilting the seat portion of the wheelchair backwards so that the user\"\"s weight is shifted away from the pressure points on the user\"\"s body. Also, the user\"\"s weight can be shifted by reclining the back frame. During the recline process it is preferable to raise the user\"\"s legs by raising the legrest. Power elevating legrests are known, and they are typically driven by linear actuators. A means for raising the legrests on a power wheelchair is disclosed in U.S. Pat. No. 5,297,021 to Koerlin et al.\nWhere a wheelchair is provided with a legrest, it is desirable to provide the legrest with a swingaway feature so that the legrests can be removed or swung out of the way to make it easier for the user to have access to or egress from the wheelchair. A particularly useful swingaway footrest is disclosed in U.S. Pat. No. 4,790,533 to Okamoto.\nOne of the problems of elevating legrests is that the legrests are typically not very substantial, giving the wheelchair user an insecure feeling. This is even a problem in power wheelchairs. It would be advantageous if there could be developed a power elevating legrest having robust characteristics. Further, it would be helpful to have a power elevating legrest with a simple to use swingaway feature to enable easy access by the wheelchair user. Another problem with existing power elevating legrests is that the linear actuators used to elevate the legrests typically end up with poorly oriented loading, resulting in badly leveraged configurations. An ideal power elevating legrest would make the maximum use of leverage for an efficient use of power and for the most compact arrangement possible.\nThe above objects as well as other objects not specifically enumerated are achieved by a wheelchair having a seat frame, legrests pivotally mounted for elevation with respect to the seat frame, and an elevation mechanism. The elevation mechanism includes a latch link having a legrest end attached to the legrest, and a pivot end. Also included is a pivot link having a frame end pivotally connected to the seat frame, and having a latch link end pivotally connected to the pivot end of the latch link. The latch link end of the pivot link and the pivot end of the latch link are pivotally connected through a latch pin. An actuator having a piston movable in forward and rearward directions with respect to the seat frame is provided to pivot the legrests relative to the seat frame, the piston having a forward end that is connected to one of the latch pin, the latch link, and the pivot link.\nAccording to this invention, there is also provided a wheelchair having a seat frame and legrests pivotally mounted for elevation with respect to the seat frame, in which the pivotal mounting of the legrests with respect to the base includes a pivot bracket that enables a legrest pivot pin to be vertically adjusted relative to the seat frame.\nVarious objects and advantages of this invention will become apparent to those skilled in the art from the following detailed description of the preferred embodiment, when read in light of the accompanying drawings."}
{"text": "Previously, in producing a high density magnetic recording tape, the magnetic layer was subjected to a calendering treatment, so as to lessen the Ra value as far as possible; (Ra value is the average surface roughness of centerline measured according to JIS 8060.) Preferably the Ra value is about 15 nm or less improves the electromagnetic conversion characteristics (Japanese Unexamined Patent Publication (Kokai) No. 5-54367). However, when the value of Ra is lessened, the coefficient of friction of the magnetic layer increases, and the trafficability of the magnetic tape is noticeably impaired. \"Trafficability\" means that coefficient of kinetic friction is small at the first pass of low speed running of a tape (10 to 1,000 cm/min.) and that adhesion of the tape at the start of running the tape is prevented.\nTo avoid the increase in friction, high molecular weight fatty acids or fatty acid esters have generally been used as lubricants. However, when fatty acids and fatty acid esters are added in small amounts, the trafficability of the magnetic tape remains poor. When these substances are added in large amounts, the strength of the magnetic layer decreases substantially, potentially lowering the durability of the magnetic tape (Japanese Unexamined Patent Publication (Kokai) No. 58-130435). While a high molecular weight fatty acid is especially effective for lessening the coefficient of kinetic friction of a magnetic tape when the tape is running at a low speed, if it is used in a large amounts, it disturbs the mutual adsorption between the ferromagnetic powder and the binder, and the strength of the magnetic layer as a coating is lowered.\nJapanese Unexamined Patent Publication (Kokai) No. 3-224123 discloses a process for the production of a magnetic tape, comprising subjecting a magnetic layer to a calendering treatment, so as to harden the magnetic layer, and thereafter subjecting the magnetic layer to contact slide action by use of a plate-like or rod-like hard metal such as carbon steel or super steel, for the purpose of reducing the abrasion characteristics and drop-out, and improving the durability of static images.\nHowever, this process has the following problems: (1) in order to form an effective lubricant film (especially a film of a fatty acid) on the surface of a magnetic layer, contact sliding procedures must be repeated several times; (2) when the contact sliding is effected by use of a metal having an Ra value higher than the roughness of the magnetic layer surface, flaws are produced such as scratches, and many errors are caused; (3) with a contact slide action in only one direction, the lubricant film of the magnetic layer becomes non-uniform or causes anisotropy, so that the coefficient of kinetic friction varies depending on the running direction (the direction in which the contact slide action was effected, and the direction contrary thereto); and (4) when put into practical use, the magnetic tape is brought into frictional contact, in the forward and reverse directions, with the magnetic head, rotary cylinder, and the like of a recording reproduction device, so that it becomes difficult to realize stable running of a tape in any of the directions."}
{"text": "A variety of techniques and apparatus have been used to satisfy the requirements of automated currency handling machines. At the upper end of sophistication in this area of technology are machines which are capable of rapidly identifying, discriminating and counting multiple currency denominations. This type of machine, hereinafter designated as a \"denomination discriminator,\" typically employs either magnetic sensing or optical sensing for identifying the denominations of bills in a stack and discriminating between different currency denominations. At a lower level of sophistication in this area are machines which are designed to rapidly count the number of currency bills in a stack, but which are not designed to identify or discriminate among multiple currency denominations. This type of machine, hereinafter designated as a \"counter,\" may include magnetic or optical sensors sufficient to enable it to discriminate between acceptable and non-acceptable bills in a stack of bills having a known denomination, but typically do not permit the machine to identify the denomination of bills or discriminate among multiple denominations of currency. Consequently, counters known in the art do not typically \"know\" what denomination they are counting until they are informed of the particular denomination by an external signal or operator.\nWhether employed in a denomination discriminator or counter, magnetic sensing is based on detecting the presence or absence of magnetic ink in portions of the printed indicia on the currency by using magnetic sensors, usually ferrite core-based sensors, and using the detected magnetic signals, after undergoing analog or digital processing, as the basis for discrimination. The more commonly used optical sensing technique, on the other hand, is based on detecting and analyzing variations in light reflectance or transmissivity characteristics occurring when a currency bill is illuminated and scanned by a strip of focused light. The subsequent currency discrimination is based on the comparison of sensed optical or magnetic characteristics with prestored parameters relating to different currency denominations, while accounting for adequate tolerances reflecting differences among bills of a given denomination. Similarly, the acceptance or rejection of a bill is based on the comparison of sensed optical or magnetic characteristics with prestored parameters defining an acceptable bill, while accounting for adequate tolerances reflecting differences among bills of a given denomination. An example of a currency handling machine using an optical scanning technique is described in U.S. Pat. No. 5,295,196, issued Mar. 15, 1994 to Raterman et al. and assigned to Cummins-Allison Corporation, incorporated herein by reference.\nCurrency handling machines (e.g. denomination discriminators or counters) known in the art typically include a system memory for storing prestored parameters associated with the magnetic or optical characteristics of the various currency denominations to be evaluated or counted. The types or denominations of currency which a machine is able to accommodate is dependent on the prestored parameters with which it has been programmed. For example, a machine designed for U.S. markets must be programmed with prestored parameters associated with magnetic or optical characteristics of U.S. currency, while a machine designed for a foreign market must be programmed with prestored parameters associated with the appropriate foreign currency. A machine designed for one market will be unable to accommodate currency from the another market unless it has been encoded with the appropriate prestored parameters for that other market. Additionally, once programmed with the appropriate prestored parameters, the system memory must be updated or supplemented periodically in order to reflect the most recent optical or magnetic characteristics of the various currency denominations to be evaluated, which may occur, for example, upon the issuance of a new series of bills.\nHeretofore, the encoding or updating of prestored parameters into the system memory of discrimination machines or counters have been accomplished externally from the machine, typically at a factory or service center. For example, in discrimination machines employing memory chips such as erasable programmable read only memorys (EPROMs), the chips are typically programmed or updated at the factory or service center and either installed in the machine at the factory or, in the case of updates, shipped to the customer for re-installation in the machine. An alternative method of encoding or updating prestored parameters may be utilized in discrimination machines employing \"flash card\" technology, such as described in U.S. patent application Ser. No. 08/715,029, now issued as U.S. Pat. No. 5,909,502, assigned to the assignee of the present invention and incorporated herein by reference. In such a \"flash card\" loading system, a flash card is programmed with the desired code and the machine may be encoded or updated by inserting the flash card into the machine, causing the system memory to become replaced with the flash card memory. Nevertheless, in either of the above prior systems, the source of the code is external to the machine, typically at the factory or service center level, and the discrimination capability of a particular machine is limited to only those bills associated with the pre-stored parameters with which it has been programmed.\nAccordingly, in view of the above-described problems, there is a need for a currency handling system that is able to accommodate currencies of several denominations and types without having been externally programmed or updated with pre-stored parameters associated with those denominations and types. Similarly, there is a need for a note counter which is able to determine the denomination of currency it is counting without having been informed of the denomination by an external signal or operator. The present invention is directed to satisfying or at least partially satisfying these needs."}
{"text": "There is a high demand for Imaging Fourier Transform Spectrometer (IFTS) hyperspectral data. Hyperspectral IFTS data is useful in a variety of fields, including atmospheric sounding (temperature, moisture, ozone), climate (greenhouse gas measurements), agriculture, and others.\nHyperspectral imaging, like other spectral imaging, collects and processes information from across the electromagnetic spectrum. The goal of hyperspectral imaging is to obtain a spectrum for each pixel in an image, such as to find objects, identify materials, and/or detect processes.\nThere are two general branches of spectral imagers: push broom scanners and related whisk broom scanners, which read images over time; and snapshot hyperspectral imaging, which uses a staring array to generate an image in an instance.\nThe human eye sees color of visible light in mostly three bands (long wavelengths, perceived as red, medium wavelengths, perceived as green, and short wavelengths, perceived as blue). Whereas spectral imaging divides the spectrum into many more bands. This technique of dividing images into bands can be extended beyond the visible. In hyperspectral imaging, the recorded spectra have fine wavelength resolution and cover a wide range of wavelengths. Hyperspectral imaging measures contiguous spectral bands, as opposed to multispectral imaging which measures spaced spectral bands.\nHyperspectral sensors collect orders of magnitude more data than traditional multi-spectral instruments.\nAn advantage of hyperspectral imaging is that, because an entire spectrum is acquired at each point, an operator needs no prior knowledge of the sample, and post processing allows all available information from the dataset to be mined. Hyperspectral imaging can also take advantage of the spatial relationships among the different spectra in a neighborhood, allowing more elaborate spectral-spatial models for a more accurate segmentation and classification of the image.\nChallenges of hyperspectral imaging include cost and complexity. Fast computers, sensitive detectors, and large data storage capacities are needed for analyzing hyperspectral data. Significant data storage capacity is necessary since hyperspectral cubes are large, multidimensional datasets. All of these factors greatly increase the cost of acquiring and processing hyperspectral data. Another challenge is finding ways to program hyperspectral satellites to sort through data on their own and transmit only the most important images, as both transmission and storage of that much data could prove difficult and costly. As a relatively new analytical technique, the full potential of hyperspectral imaging has not yet been realized."}
{"text": "The present invention relates, in general, to electronic devices. More particularly, the present invention provides a transparent solar cell and method of its manufacture.\nSolar energy provides many advantages over traditional energy sources. For example, energy from the sun is virtually unlimited and easily accessible throughout the world. It does not require the extraction of a natural resource from the ground to obtain the energy and it can be converted to electricity in a manner that is not harmful to the environment. Solar energy is available whenever the sun is shining and can be collected and stored for use when no light source is available. Therefore, if it can be harnessed economically, it provides an environmentally friendly source of energy that does not deplete or destroy precious natural resources. This is in stark contrast to the use of fossil fuels that are of limited supply and which cause environmental damage with both their use and extraction processes. The use of fossil fuel also requires a constant source of raw materials that may be difficult obtain in many circumstances.\nMany different applications benefit greatly from the use of solar energy. For example, buildings and automobiles, with their broad surfaces that are exposed to the sun's energy for much of the day, can use that energy to provide some or all of their energy needs. Various solar cells have been developed using different fabrication techniques to take advantage of this energy source.\nOne type of solar cell is formed with crystalline silicon. For these solar cells, crystalline silicon is formed by melting silicon and drawing an ingot of crystalline silicon of the size desired. Alternatively, a ribbon of crystalline silicon can be pulled from molton silicon to form a crystalline silicon solar cell. A conductor is placed on either side of the crystalline silicon to form the solar cell. These processes use high temperatures and the solar cells are expensive to manufacture. Packaging is also difficult and expensive and creates a rigid structure. Their maximum size is limited by the manufacturing process. It is difficult to slice the resulting crystalline silicon thin enough to provide a transparent or flexible solar cell. However, these structures are very efficient (relative to other types of presently available commercial solar cells). As such, crystalline solar cells are used primarily for applications where efficiency is more important than cost and where the structures do not need to be flexible. For example, these are commonly used on satellites.\nAnother type of solar cell is formed with polycrystalline silicon. These may be formed as thin layers on wafers and can thus be made thinner than crystalline silicon solar cells. As is well known in the art, polycrystalline silicon can be formed by heating amorphous silicon. Typically, amorphous silicon begins to crystallize at temperatures greater than about 1400.degree. C. Because of these high temperatures, known processes can only use substrates with high melting points. These processes are not appropriate for substrates made of plastics or other materials that melt at lower temperatures. In the manufacture of flat panel displays, it is known to use lasers to form polycrystalline silicon thin film transistors (TFTs). Such use has not included the formation of P-N junctions or solar cells which presents its own set of challenges. Moreover, these manufacturing processes generally formed single transistors and were not used to form large sheets or areas of polycrystalline silicon. Further, lasers have been used in the manufacture of solar cells, but only as a tool to mechanically form (slice, pattern, etch, etc.) the solar cells.\nAnother type of solar cell has been formed using doped layers of amorphous silicon. These are not subject to some of the problems inherent in the previously described crystalline silicon or polycrystalline solar cells. First, amorphous silicon can be formed using low temperature processes. Thus, it can be formed on plastic and other flexible substrates. They can also be formed over large surfaces. Second, the processing techniques are less expensive. Nevertheless, amorphous solar cells introduce other significant limitations not found in crystalline silicon or polycrystalline silicon solar cells. For example, hydrogen is generally added during the manufacturing to increase the efficiency of the cell. Amorphous silicon solar cells tend however to lose this hydrogen over time, causing reduced efficiency and reduced usable life. Moreover, amorphous silicon solar cells are not transparent. Thus, they are not appropriate for many applications. For example, buildings and cars with solar cells can be unsightly, and the solar panels may block the view of the outdoors or access to outside light indoors. Also, portable electronics often place a premium on size and surface area. Some devices have displays that cover most--if not all--of the exposed surface of the device. Therefore, it is often undesirable or impossible to mount a traditional amorphous silicon solar cell on the device.\nAttempts have been made to solve this transparency problem by making transparent panels from existing solar cell processes. One method has been to take advantage of the \"window shade effect\" whereby solar cells are formed on a transparent substrate with gaps between adjacent solar cells. This allows some light to pass through to create a transparent effect. The larger the gaps, the more transparency the device has. A disadvantage of this technique is that much of the space is unused, therefore the efficiency of the device is less than it would be if all of the surface area were used for solar cells. Of course, devices of this type also suffer from the problems inherent to the type of cell used. For example, if based on amorphous silicon, these devices suffer from the hydrogen loss exhibited in other amorphous silicon devices.\nOther work has been done at making transparent solar cells using materials other than silicon (for example, cadmium telluride (CdTe)). These cells suffer from the challenges inherit to using materials other than silicon.\nThus, a new solar cell and method of fabrication that will avoid these problems is desirable."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the management of context information, and more particularly to a scalable and extensible system and method for collecting, maintaining, and disseminating such information.\n2. Description of the Related Art\nWith the proliferation of small computing devices, wireless networking, and sensor technology, comes the ability to collect and maintain a wide variety of information about people, places, and objects. This information offers the opportunity to make application programs aware of the context of those people, places, and objects, which in turn will reduce the required amount of human attention to application programs. By tracking and interpreting the context of a person, place or object, context systems can automatically configure computing devices and direct information transfer without user intervention.\nOne system for managing context information (such information sometimes referred to in the literature as “awareness information”) is described in commonly assigned and co-pending U.S. application Ser. No. 09/511,977, filed Feb. 24, 2000 (the “prior application”), incorporated herein by reference. The prior application describes a system for collecting context information and storing such information on an awareness server that can be contacted by a closed set of users and devices. Information is provided to the awareness server from registered users and machines by means of information packets that include a record identifier and an encrypted form of the data. When the awareness server receives an information packet, the identifier is checked for validity, and if valid, the attached context information is decrypted, further data source validity checking is performed, and the preferences of the information source and object are checked to determine authorized handling of the information. While the system described in the prior application is configured to receive and maintain context information of varying types, there is no means described for easily extending the system, after initial setup, to handle more diverse context information. In addition, the prior application does not provide a means for scaling to handle substantially greater volumes of context information of a given type.\nAs context computing becomes more prevalent, so will the need for simple communication between users and devices and the context service to which they subscribe. In the past, in order to update a context system to handle additional types of information, it was necessary for a programmer to manually alter the server code. With large context systems, this can become a bottleneck to system upgrades.\nClearly, as computing systems become more ubiquitous, there will be a need for the ability to quickly and easily extend context systems to handle information of diverse types, as well as to scale to the capacity of such systems to meet ever larger user communities. In addition, an easier method of requesting and exchanging information with context services will also be necessary."}
{"text": "The endocrine pancreas consists primarily of islet cells that synthesize and secrete the peptide hormone glucagon, insulin, somatostatin and pancreatic polypeptide. Insulin gene expression is restricted to pancreatic islet β-cells of the mammalian pancreas through control mechanisms mediated in part by specific transcription factors. In other cells the insulin, other pancreatic hormones and specific peptidases genes are trancriptionally silent. The homeodomain protein PDX-1 (Pancreatic and Duodenal Homeobox gene-1, also known as IDX-1, IPF-1, STF-1 or IUF-1) plays a central role in regulating pancreatic islet development and function. PDX-1 is either directly or indirectly involved in islet-cell-specific expression of various genes such as for example insulin, glucagon somatostatin, proinsulin convertase 1/3 (PC 1/3), GLUT-2 and glucokinase. Additionally, PDX-1 mediates insulin gene transcription in response to glucose."}
{"text": "Commonly, fastener-driving tools, such as powder-actuated tools, are arranged to drive fasteners of a known type comprising a shank defining an axis and having a tip at one end, a head integral with the other end of the shank, and a washer carried by the shank with an interference fit. Such fasteners are exemplified in Almeras et al. U.S. Pat. No. 4,824,003.\nIn such a fastener, the washer is carried near but in spaced relation to the tip and is moveable axially toward the head when the fastener is driven with the washer bearing against a workpiece. The head diameter and the washer diameter are approximately equal.\nAs exemplified in Almeras et al. U.S. Pat. No. 4,824,003, it is known for such a tool to be muzzleloaded with such fasteners, which are loaded one at a time. As exemplified in Pfister U.S. Pat. No. 4,881,643, it is known to load a plurality of different fasteners into a powder-actuated tool, via a carrier strip fed laterally into the tool.\nA common use of a powder-actuated tool, as exemplified in Almeras et al. U.S. Pat. No. 4,824,003, is to attach metal decking members to steel structural members or concrete floors. For such a use, it would be highly desirable to adapt such a tool so as to facilitate its use by a standing worker. Neither a muzzle-loaded tool nor a strip-loaded tool would be entirely satisfactory, since the worker would have to lift the tool or to stoop whenever it was necessary to reload the tool.\nThus, there has been a need, to which this invention is addressed, for a better approach to loading fasteners into a fastener-driving tool, such as a powder-actuated tool, so as to facilitate its use by a standing worker."}
{"text": "The deposition of thin metallic films on substrates by the decomposition of metallo-organic solutions is a technology which has been known for many years. Although the technique has been known for a long time, its primary application has been in the decorative trade, with the main emphasis on producing films of precious metals on ceramic and glass articles. The earliest known application was bright gold for decorating porcelain and this process had been used since 1830.\nAlmost all of the work prior to 1980 used metallo-organic compounds derived from resins or other natural products and were commonly called metal resinates. These resinates were and still remain suitable for most applications in the decorative trade, but the variability of chemistry, which is an inescapable result of the preparation from natural products, was a major impediment to the extensive development of technical applications of the metallo-organic decomposition (MOD) process. It has only been in recent years that MOD films have been produced from pure, well characterized compounds.\nRecent applications of MOD processing are in processing technical ceramics where most of the emphasis is on fine particle processing. Sol-gel methods use metal alkoxides as the starting materials and have general formula M(OR).sub.n where M and R are metals and alkyl groups, respectively. Metal alkoxides are hydrolyzed easily and yield oxides, hydroxides, and hydrated oxides in crystalline and amorphorous forms. The fine powders derived from the metal alkoxides have high surface reactivity, making it possible to use low sintering temperatures and obtain uniform grain size bodies with the desired electrical properties. Others have prepared films from the metal alkoxide solutions by dipping substrates into the solution. The barium titanate films prepared by this method have had very good electrical properties. The metal alkoxide method, however, is a sol-gel method and involves several lengthy steps in processing.\nThe metallo-organic decomposition technology investigated by R. W. Vest, G. M. Vest and others produce the metal oxide films or powders directly and more quickly than those produced form the sol-gel methods. For example, silver films were fabricated by thermally decomposing the MOD silver inks on silicon for use as collector grids for photovoltaic cells. These collector grids required good adhesion, low contact resistance, low sheet resistance, and long term stability. Also, their deposition methods should not degrade the n-p junction. The fired films were 100% solder leach resistant, had good line definition and excellent long term adhesion. The resistivity was very close to that of bulk silver, indicating the films were very dense. This method of metallizing solar cells was also very inexpensive as compared to sputtering or vapor depositing on Si.\nThere are many advantages of MOD processing compared to alternate techniques for producing metal and ceramic films. The MOD process yields the equilibrium phases of the desired systems at relatively low temperatures, which circumvents the problem of selective volatility of different species. In general, the low temperature processing yields extremely find grain size polycrystalline films; in many cases the initial inorganic films are amorphous to X-rays. This allows for precise control of grain size by annealing after preparation of the films. The low temperature processing and the achievement of equilibrium phases is primarily due to the fact that the formulation that is deposited on the substrate is a true solution, and hence, the mixing of the various ingredients is on the molecular level. This ultimate mixing and high reactivity also can be used to an advantage in preparing very dense films. In most cases, the films with near theoretical density can be achieved. Starting from solution also leads to films with extremely uniform composition over large areas, and allows for uniform doping in the ppm range. High purity can be maintained during MOD processing by appropriate care in the various processing steps.\nThere are some intrinsic limitations to MOD processing, however. The volume change in going from the deposited wet film to the fired inorganic films is always large. In addition to requiring care during thermal processing, this large volume change means that the fired films will always be thin. This limitation of film thickness can be overcome by the multilayer approach. A second intrinsic limitation is one of the advantages cited in the previous paragraph, namely, that thermodynamic equilibrium is achieved very rapidly because of the extremely high reactivity upon thermal decomposition. Many of the electronic films in use today have their desirable properties because of the non-equilibrium microstructure. These non-equilibrium microstructures cannot be duplicated by MOD technology but alternate approaches to achieve equivalent electrical properties with an equilibrium microstructure can often be utilized. Another limitation due to the thermodynamic equilibrium achieved is that only oxides or only metals of certain elements can be produced."}
{"text": "THIS invention relates to a method and system for purifying or cleansing a gas stream or gaseous body containing one or more contaminants.\nIt is well established that outdoor air pollution can cause serious human illness. Air pollution can cause burning eyes and noses, itchy irritated throats as well as breathing difficulties. Some chemicals can cause cancer, birth defects, brain and nerve damage, long-term injury to lungs and even premature death. As a result, protection of the public health is a driving force behind the laws and regulations that are being, and have been, put into place to clean up the environment. However, ridding outdoor air toxic contaminants has often overlooked indoor air quality (IAQ).\nThe results of numerous scientific investigations have revealed many reported outbreaks of what is known as Sick Building Syndrome (SBS) and Building Related Illness (BRI) as a direct result of indoor air pollution (IAP) in office and other public access buildings.\nSBS is used to describe the situation in which building occupants experience acute health and comfort defects that are linked to time spent in the building. BRI is the term used when diagnosable illness and symptoms are identified and directly attributed to airborne building material contaminants.\nIn modern society, it is believed that most people spend as much as 90% of their time indoors. The US Environmental Protection Agency (EPA) has indicated that IAP may be between 2–5 times higher than accepted levels, and up to 100 times higher than outdoor pollution levels.\nThe World Health Organisation (WHO) Committee Report on IAQ suggested that 30% of new and re-modelled buildings globally are subject to excessive complaints of IAQ and have unusually high rates of SBS and/or BRI. For example, according to the SA Singapore embassy, 92% of all absenteeism in Singapore can be attributed to SBS and/or BRI.\nIAQ is rated within the top 5 of the US EPA's most critical global concerns alongside ozone depletion, global warming and the like. It has been reported, for instance, that in 1989, the US economy suffered a staggering US $65 bn loss due to SBS and BRI. The absenteeism levels, due to poor health caused by IAP, resulted In this huge productivity decrease."}
{"text": "Pipeline pigs are used to clean pipelines. Pipelines run the hazard of forming accumulations or deposits on the interior. Such deposits accumulate over a period of time from the product flowing through the pipeline. In some instances, this is no problem whatsoever. The deposit may form very slowly and, even after formation, may be removed easily because it is relatively soft. In other circumstances, the deposit placed on the interior of the pipeline may be very hard. Depending on the locale, artesian water will include a significant mineral content. The formation of internal pipe scale is accelerated where the water is hot. Where highly mineralized artesian water is heated and flows through a pipeline, scale coating in the pipe is accelerated. It is possible to coat a pipe on the interior to almost totally plug the pipeline. The internally coated pipe scale must be cleaned. Sometimes, cleaning can be achieved with other types of pigs. However, it has been found that the pig described in this specification is superior in many cleaning operations. It is particularly able to break up the scale on the interior of a pipe so that it can be removed in large pieces.\nPigs equipped with thousands of small points have been known before, and they will eventually remove pipe scale from a pipe. One such pig is supplied by the firm Polly-Pig by Knapp, Inc., of Houston, Tex., and is sold under the trademark \"Superjavalina.\" It is a very good pig for many purposes. One of its virtues arises from the incorporation of thousands of small, relatively short metal chisels which cut the pipe scale away. So to speak, they collectively cut away a cloud of dust. By contrast, the pig of this disclosure cuts the scale only in a few places, but it tends to cut deeply and breaks away large pieces or chunks of the pipe scale. As the large pieces are broken from the unwanted pipe coating, they are flushed by the fluid flow which propels the pig through the pipeline. This enables the line to be cleaned and flushed with just a few passes of the pig. This enables the pipeline to be restored to full service more promptly.\nOne advantage of the highly inventive pig of this disclosure is that it appears to work even better with coatings which are relatively harder than most coatings. Indeed, the relative hardness and brittleness of the coating enhances the operation of the pipe cleaning pig of this disclosure. Enhancement is obtained in that soft, pliant coatings will permit the sharp, pointed stud to dig and gouge, but will not necessarily initiate cracks running through the coating which break the coating free. By contrast, a brittle pipe coating which is cut in a vigorous fashion by the chiseling point of the stud associated with this pig tends to break loose in large pieces as cracks and fissures are propogated from the line of contact of the stud with the pipe scale. This breaks up the pipe coating and permits it to flake away in larger pieces. It has been suprisingly learned that it is, therefore, better to reduce the number of metal chisels striking against the internal scale. This is borne out by the reduction in points of contact against the internal scale from the \"Superjavalina\" pig which approximates hundreds, or perhaps thousands, of points per square inch to the present disclosure which utilizes only a few metal chisels on the entire pig. This reduction of several hundredfold to a few stud points on the body of the pig enables the hard internal pipe scale to break up in large pieces and permits its removal more rapidly.\nThe present apparatus has the further advantage of removing other kinds of hard coatings. For instance, highly irregular coatings are easily removed. One kind of highly irregular coating is that found in pipelines which carry seawater and which do not have sophisticated water filtration apparatus. Such pipelines (for instance, communicating with a cooling tank) often accumulate aquatic shell deposits of a very random distribution, thickness and hardness.\nThe improved pig of the present apparatus undergoes significant stress in operation. The tip of the stud may undergo significant bending as it impacts a brick-like pipe scale deposit. The brick-like deposit may break and shatter, but only if the stud maintains its relative position in the pig as the pig traverses the pipeline. It is important that the stud be firmly anchored in the pig to resist the wear and tear of operation. Failure of the pig can be manifested by tearing the stud out of the pig body. One important advantage of an embodiment disclosed herein is the incorporation of an enlarged shoulder or facing area near the bottom end of the stud which anchors the stud in position and prevents its tearing free. Thus, the tip of the stud may be violently deflected, but the resilience of the body of the pig enables the stud to accommodate such deflections while the stud remains intact with the body. The transverse shoulder anchors the stud in position.\nThe present invention has as one of its objects the provision of a pipeline pig which is pressure propelled through a pipeline with a plurality of studs mounted in spaced locations and including spirals thereon which tend to rotate the pig as it traverses a pipeline, the studs gouging and chiseling away hard pipe scale."}
{"text": "1. Field of the Invention\nThe present invention relates to image data filtering.\n2. Background Art\nAn increasing number of devices are being produced that are enabled to capture and display images. For example, mobile devices, such as cell phones, are increasingly being equipped with cameras to capture images, including still snapshots and motion video images. Images captured by such devices can frequently be viewed on displays of the devices, as well as being transferred from the devices for viewing elsewhere. To view the images on relatively small devices, the images typically must be viewed on small display screens that are not capable of viewing the full resolution of the captured images. Thus, such devices must include at least limited image processing capability to down-size the images for viewing on the small display screens.\nMany mobile devices have limited processing capability due to cost, power consumption, and size constraints. However, the processing of captured images, especially the processing of video, is very computationally intensive. For example, many mobile devices have cameras capable of capturing images of 2 MegaPixels (MPel) or more. Thus, a processor of such a mobile device must be capable of processing a large amount of data for each captured image. Furthermore, encoding and decoding (e.g., QCIF) of image data may need to be performed by the processor at frame rates such as 15 fps and 30 fps, respectively, as well as the performance of other functions.\nTo deal with such high-resource demanding tasks, mobile device developers have resorted to including high-powered processing chips in the devices, which have higher clock rates and larger on-chip cache memory, or to including dedicated video/image processing chips. Such approaches result in higher cost devices with higher levels of power consumption, which may not be desirable in battery powered mobile devices.\nThus, ways are desired for handling an image processing workload in devices, such as mobile devices, without significantly raising device costs and power consumption levels."}
{"text": "Electronic pens have been developed as a complement to keyboards and mice for controlling computer-based systems and inputting information to them. Such a pen is associated with an internal and/or external position detector which electronically records the pattern of movements of the pen over a base and generates an electronic copy thereof, below referred to as electronic handwriting.\nThe position detector can be based on, for instance, recording of a position-coding pattern on the base, mechanical detection of movements, use of accelerometers and/or gyroscopes, triangulation of electromagnetic or acoustic signals, pressure detection in the base, recording of changes in an electromagnetic field, or a combination of such techniques.\nIn many electronic pens, the position detector is activated based on an output signal from a pressure sensor which is mechanically connected to the pen stem to detect when the pen is made to abut against a base.\nWhen reproducing electronic handwriting on a display or a base, readability is frequently dependent on the fact that the reproduction also reflects dynamic parameters of the handwriting, such as the force of application of the pen stem during writing. The above pressure sensor should thus be designed to emit an output signal that is representative of the force of application.\nAn electronic pen of this type is known from e.g. U.S. Pat. No. 6,252,182 B1. The pen has a pen stem whose inner end is connected with a control block which is slidingly received in a cylindrical cavity in the body of the pen. Opposite to the control block, there is arranged in the cavity an electrode plate with an inner circular conductor and an outer helical conductor. A helical spring is arranged round the control block in the cavity for biasing the same away from the electrode element into abutment against a stop lug. On its side facing the electrode element, the control block has a pointed application element of conductive rubber. In an unloaded initial position, the application element with its point is received in a central guide hole in the electrode element and is in electrical contact with the circular conductor. When exerting a load on the pen stem, the application element is gradually urged against the electrode plate so that the contact surface between them increases, whereby the electrical resistance over the electrode plate gradually changes. Based on the voltage across the electrode plate, a measure of the pen stem's application force on a base is thus obtained.\nA pen of the above construction, however, suffers from several drawbacks.\nThe sensor arrangement of the pen is relatively complicated to assemble, which makes production as well as repair and maintenance difficult and expensive.\nMoreover the sensor arrangement of the pen will probably have a limited life by the point of the application element being quite heavily worn when cooperating with the guide hole in the electrode plate.\nA further drawback is that the characteristics of the sensor arrangement, i.e. its functional relationship between application force and output signal, can only be modified within narrow limits and then by way of radical construction changes of the sensor arrangement.\nFurther examples of pressure-sensitive electronic pens are given in U.S. Pat. Nos. 6,104,388, 5,895,895, 5,571,997 and U.S. Pat. No. 5,548,092."}
{"text": "1. Field of the Invention\nThe pesent invention relates to a supervisory device for wrong operation of varying quantity measuring apparatus, and more particularly to a supervisory device for wrong operation of varying quantity measuring apparatus which measures a physical varying quantity by relative movement between a main scale and an index scale.\n2. Description of the Prior Art\nThere has been a well-known physical varying quantity measuring apparatus which measures a certain physical varying quantity by relative movement between a main scale and an index scale, and is of wide use for the measurement of length and physical varying quantity in length measuring instruments, coordinate measuring machines, etc. or in positioning or the like of machine tools, or for the other purposes.\nThis varying quantity measuring apparatus includes a main scale and an index scale which are moved relative to each other in accordance with a physical varying quantity and an encoder which detects the relative movement of such two scales and measures the requested physical varying quantity based upon the relative movement of such two scales which can be detected by this encoder.\nThe encoder mentioned above detects the relative movement of both two scales as an electric signal by optical, electromagnetic or well-known means using contact resistance, and further converts the detected signal into a count pulse signal in accordance with the relative movement per unit time to supply this to a counter provided in the encoder. The counter counts the supplied count pulse signal, and the counted value of the count pulse signal is displayed as a digital value representing a measured value of the physical varying quantity or is used for an other control signal.\nAccordingly, this varying quantity measuring apparatus can measure the physical varying quantity with accuracy as a predetermined digital signal by means of relatively moving the main scale and the index scale in accordance with the desired physical varying quantity.\nThe varying quantity measuring apparatus mentioned hereinabove, however, clears the counted value of the count pulse signal stored in the encoder, when the electric power stops or lowers the voltage to great extent for some reason, and restart counting the count pulse signal from the very beginning, when the electric voltage recovers up to the optimum level. Accordingly, when the electric power stops or lowers its voltage, there arise measuring errors between the relative moving quantities of the main scale and the index scale which are converted into the electric signal by the encoder to be displayed and the actual moving quantities of these two scales. Furthermore, it is a problem that the error occurrence is not known to those at measuring work in most cases.\nEspecially, in such a case that the electric power suspension or low voltage is instantly caused by the start or the stop of a large machine, the sudden load change, or the like for example, the encoder immediately restart counting the count pulse right after the counted value of the count pulse signal is cleared, which is not known to those at measuring work in most cases, and an effective countermeasure has been requested."}
{"text": "Generating a session key from a known shared master key is known in the art. However, previous approaches have a number of limitations. For example, storing a large number of keys for session communication has an impact over the performance of a system and may also make the system vulnerable if the security of the system is jeopardized. Another way of creating a session key is by using one-way hash functions. However, generating a session key using a one-way hash function requires complex mathematical operations that come at the expense of computation time. Yet another method of creating a session key known in the art is by exchanging a random signal between each of the communicating users and then using the exchanged random signals to form the new key. However, this method creates overhead, as both parties or communicators are required to compute random numbers for generating the key. Accordingly, there is a need for a method and system of generating a session key so that both the computing time and the memory requirement can be reduced."}
{"text": "Some types of conventional computer systems are used to manage hardware resources to perform particular tasks. Example hardware resources include stimulus resources and measurement resources. Stimulus resources apply stimuli to one or more devices. For example, a stimulus resource can apply a voltage or a current to a device, such as a circuit, switch or relay. In turn, the device can direct a voltage or a current to a particular site, or direct pneumatic or hydraulic fluid to move mechanical arms. Measurement resources sense or measure a response from a device. For example, a measurement resource can include a sensing circuit that provides precise voltage or current measurement values, or bits of information that indicate the location of an arm.\nGenerally, the computer system sets stimulus resource values prior to calling a measurement resource. Often, the computer system must delay performing the measurement operation while waiting for the stimulus to settle to a stable value or state. This delay time is referred to as a settling time delay. Multiple stimulus resources introduce multiple delay values that must be managed to optimize the time it takes to perform a particular task or series of tasks.\nIn one type of computer system, referred to as a circuit test system, stimulus resources apply stimuli to an integrated circuit and measurement resources measure the circuit's response. Stimulus resources provide stimuli in the form of logic levels or analog voltages and currents. Measurement resources take measurements of logic levels, analog voltages and currents, and timing values.\nCircuit test systems perform tests on a circuit, referred to as a device under test (DUT). Example types of tests performed by a circuit test system include DC parametric tests, functional tests, and AC propagation delay tests. Parametric tests check DC voltage and current handling characteristics of inputs and outputs. Functional tests check DUT functionality using test vectors, which preferably exercise all functional aspects of the DUT. Propagation delay tests check the time it takes for a signal to propagate through the DUT, referred to as the AC characteristics. Often, each test type is performed multiple times under different conditions. Hardware resources are typically changed for each test.\nA user, such as a test engineer, typically supplies a test program to test a DUT in a particular test sequence. The test program includes instructions that change hardware resources to accomplish testing. This test program is loaded onto the test system and executed to alter hardware resources and test a DUT. In one example sequence, the test program first supplies instructions to power the DUT at a selected source voltage. Source current is measured and checked against limits to detect shorts and opens on the power line. After a DUT has passed this test, DC parametric tests are performed to check the integrity of each input and output. The DC parametric tests are performed by applying selected voltages and currents to particular input and output pins. Test vectors are often employed to obtain proper testing states at the output pins. The test system measures current and voltage values from selected pins.\nAfter DC parametric tests, functional tests are performed. Functional tests can include stepping through a set of test vectors where selected logic levels are applied to input pins, and output pins are checked against expected logic states. A functional test can be performed multiple times with different voltages on the power pin. Next, propagation delay tests are performed to ensure the DUT operates up to speed. Often, when hardware resources are altered, the test system must wait for the voltage or current to settle before taking a measurement. Waiting for the voltage or current to settle helps to ensure an accurate measurement, but increases test time and raises device cost.\nTest engineers write test programs to optimize the balance between testing accuracy and speed. Test system manufacturers publish tables of hardware resource settling time delay values. These settling time delay values are usually worst case settling times for hardware resources. In one optimization method, test engineers write test programs to provide a delay after each resource change, which results in long test times. In order to reduce these test times, test engineers typically measure settling times and use measured delay values instead of published delay values. Using the measured delay values ensures accurate measurements and reduces test time.\nIn another method, test engineers change multiple resources before providing one delay value for all changed resource settling times. In this method, test engineers use the longest settling time for all changed resources. This settling time is taken from a published table or, to reduce test time further, the test engineer measures settling times for all changed resources and chooses the longest measured settling time. However, measuring settling times can be a time consuming and frustrating task, especially, if settling times change from one test system to another. In each of these methods, test program instructions are not executed while the resources settle. Also, the burden of optimizing the test program is placed on the test engineer.\nIn an effort to reduce the burden on the test engineer, test system manufacturers typically include a settling time delay at the end of each function that changes a hardware resource. However, this causes the test system to incur a settling time delay for each hardware resource change. To allow for multiple resource changes before a single delay, test system manufacturers provide test engineers with a switch to shut off the delay for a particular resource change. Using the switch, the test engineer can set one delay for a group of resource changes. However, the burden is still on the test engineer to figure out which delay or group of delays to use. Also, the test system cannot process instructions while the hardware resources settle.\nFor reasons stated above and for other reasons presented in the description of the preferred embodiments section of the present specification, an improved delay management system is desired that does not have the above problems associated with using published tables of settling time delay values, measuring settling times, and/or placing the optimization burden on the test engineer."}
{"text": "Electric or hydraulic actuators used on trailers often require high power drive. Applications such as dump trailers, liftgates (on trailers or motorized vehicles), heavy load leveling systems and similar functions are frequently implemented on vehicles not equipped with high powered hydraulic or electric distribution systems. High power actuators for mobile systems typically require high intensity current pulses in the order of hundreds of amperes at low voltage such as ˜12V or 24V, to be sustained for periods in the order of the minute. Such power and energy usually exceeds the capacity of the inherent electrical distribution system and require the use of an auxiliary battery. Battery technologies used in such applications include Lead-acid and Lithium-ion.\nLead-acid batteries are frequently used because of their low cost and availability but can be considered relatively heavy/bulky and require long charging times. They provide a limited number of cycles through their lifetime when used in harsh environment with frequent charge/discharge cycles at high amplitude current pulses. They require maintenance and regular replacement. When left discharged and exposed to cold temperatures, they can freeze and suffer permanent internal damage leading to battery failure.\nLithium-ion batteries, when properly sized, are capable of handling high intensity and long duration current pulses. However, they typically require sophisticated battery chargers for reliable and long term use. Moreover, during long term storage of the system, Lithium-ion batteries must be cycled and recharged periodically to avoid internal damage and suffer degradation causing loss of performance, making them unsuitable for certain applications.\nAccordingly, while the prior technology was satisfactory to a certain degree, there always remains room for improvement."}
{"text": "Collapsible kayaks having several designs have been commercially available for many years. Some of the better known types are assembled from a plurality of disconnected stringers and formers which are joined together with latches and clamps to form a relatively loose framework, in which the bow and stern halves are temporarily hinged together. A skin or liner is stretched over the framework having bow and stern-shaped pockets and a central upper opening adjacent the cockpit structure. The opening receives the two halves of the framework, sometimes separately, which halves are then rigidly connected together inside the skin. In some designs an over-centering means or an equivalent stretches the skin to form a relatively taut skin around the framework. Commonly the skin is strained somewhat to provide a reaction to force on the framework and thus is an integral structural part of the kayak, that is, without the skin the kayak frame has insufficient rigidity.\nSuch designs, when disassembled, result in a large number of disconnected pieces and require a considerable time to re-assemble, and require skill and care in following the instructions. It is not unusual to accidently lose an important piece when the kayak is disassembled. If the kayak is to be transported in its collapsed state at least two bags are required, one of the bags containing the stringers, formers, fasteners, etc. of the framework, and the other bag containing the skin. Whilst such a kayak has the advantage of being collapsible for fitting in a confined space, these structures are usually not designed primarily for easy carrying on a person's back.\nDue to these difficulties, there is a significant demand for a watercraft that can be collapsed into small packages for storage or transport. The simplest of these are inflatable boats, which use a system of air-filled pneumatic tubes to provide form, structure and flotation. While such watercraft are compact and inexpensive, they are typically slow and hard to maneuver, due to shape limitations imposed by pneumatic systems. Thus, there is a need for a kayak that is more easily transportable than a traditional rigid-bodied kayak, yet more durable and reliable than an inflatable kayak."}
{"text": "There is conventionally known, as a vehicle seat, a vehicle seat having a cushion material formed from polyurethane foam (hereinafter called urethane). As the cushion structure of such a vehicle seat, a structure is widely used in which a cushion material made of urethane is placed on a plate or a spring material, such as a contour mat (trade name) or the like provided at a seat frame or a seat back frame, and these are wrapped in a fabric material.\nTherefore, it is known that the shape (the design shape) and the elastic characteristic of the cushion material have a great effect on the body pressure dispersability of the seated person and the vibration absorbability. Further, by structuring the cushion material by layering urethanes which have various characteristics, a cushion material can be obtained which has a spring characteristic (elastic characteristic) which is near the spring characteristic of the muscles of the buttocks or the like of the seated person. However, in such a structure, there are the problems that there is the sensation that the restoring force is insufficient and that the weight is heavy.\nThus, a seat has been conceived of in which, as the cushion material taking the place of urethane, a cushion structure is structured by stretching, over a seat frame, a pair of ground knit fabrics and a two-dimensional knit fabric or a three-dimensional solid knit fabric formed by connecting threads which are disposed between the ground fabrics (see, for example, Japanese Patent No. 5013089). The cushion material formed from this three-dimensional solid knit fabric or two-dimensional fabric is an elastic structure which is difficult to weaken, and is thinner than urethane, and exhibits an elastic characteristic in place of urethane.\nHowever, in a conventional vehicle seat using such a three-dimensional solid knit fabric or two-dimensional fabric as the cushion material, the two-dimensional knit fabric is stretched over the seat frame at a tension such that the elongation is in the range of from 5% to 20%, and the three-dimensional solid knit fabric is stretched at a tension such that the elongation is less than 5%. Therefore, when a person sits down, large forces arise due to tension at the portions where the convex portions of the human body, such as the ischia, the coccyx, the shoulder blades and the like, contact the cushion material (clothes are ignored). Therefore, this becomes a cause of the supporting pressure becoming strong and a sensation that a foreign object exists arising. Further, if this acts over a long time period, there is the problem that pain arises in the muscles at the peripheries of the aforementioned ischia, coccyx, shoulder blades and the like. In particular, the muscles and the blood vessels in the peripheries beneath the ischial tuberosities, which support the majority of the body weight of the human body, are compressed, and there are cases in which pain or numbness due to interruption in blood circulation arises.\nFurther, in the above-described conventional seat, the two-dimensional knit fabric is stretched directly over the frame. Therefore, the problem is known that it is easy for there to arise a phenomenon, which is the so-called hammock sensation phenomenon, in which it is easy for the seated person to move in the front-rear and left-right directions within the frame, and due to the input of vibration at the time when the vehicle travels for example, the seated posture of the seated person changes or the seated person slides forward on the seat."}
{"text": "Aprotinin is a polypeptide of 58 amino acids with a molecular weight of 6512 daltons. This single chain protein has three disulfide bridges, linking cysteines at position (5-55), (14-38), and (30-51). These disulfide bonds, along with numerous hydrogen bonds and hydrophilic interactions, maintain the compact acid-stable tertiary structure of this molecule. Aprotinin has been well studied because of its small size, tremendous stability and function.\nAprotinin is a potent, naturally occurring enzyme inhibitor which inhibits proteolytic enzymes of the serine protease family. It is commonly isolated from a variety of bovine tissues including lung, parotid gland, spleen, liver, and pancreas. Specific enzymes known to be inhibited by aprotinin include trypsin, chymotrypsin, kallikrein, and plasmin. Aprotinin is also known as Bovine Pancreatic Trypsin Inhibitor (BPTI) and Trypsin Kallikrein Inhibitor (TKI).\nAprotinin is the active ingredient of a drug which has found wide use for the treatment of traumatic hemorrhagic shock, inhibition of plasmin, preservation of blood cell function in acid citrate dextrose (ACD)-stored blood as well as many other applications. In general, aprotinin is useful in reducing pathologically elevated proteolytic enzyme activity to normal levels.\nAprotinin is also useful in preventing proteolytic damage of recombinant products, such as proteins, polypeptides, enzymes, proenzymes, and preprotein, during purification and isolation procedures. The optimum amount of aprotinin for use in total inhibition of proteinases must be determined on a case-by-case basis; however, micromolar concentrations of aprotinin are usually sufficient.\nAprotinin is used extensively in radioimmunoassays to prevent proteolytic damage of the ligand. Some of the radioimmunoassays that use aprotinin as a proteolytic enzyme inhibitor are assays of adrenocorticotropin, calcitonin, beta endorphin, glucagon, and renin somatostatin.\nConsidering the many clinical and research applications of aprotinin, a ready supply of aprotinin, such as that which would result from fermentation of aprotinin-expressing recombinant organisms, will be of great value to the medical and biotechnology fields.\nSince isolation from natural sources is technically difficult, expensive, and time consuming, recent efforts have centered on the development of efficient recombinant methods for the production of APR.\nOf the hosts widely used for the production of heterologous proteins, probably E. coli and Saccharomyces cerevisiae (Baker's yeast) are the best understood. However, E. coli tends to produce disulfide-bonded proteins such as APR in their reduced forms which frequently are not stable in the presence of endogenous bacterial proteases, and which tend to aggregate into inactive complexes. Attempts to overcome this problem, e.g., by employing a suitable leader sequence in order to produce soluble APR which could be readily recovered from the cell broth, resulted in other inconveniences, especially during purification of the product, since the bulk of the desired protein was associated with the cell paste.\nFor example, Auerswald, Schroder and Kottich, in Biol. Chem. Hoppe Seyler (Germany, West) 368, 1413-1425 (1987) disclose the synthesis, cloning and expression of recombinant aprotinin in E. coli using a synthetic gene formed by the fusion of DNA fragments of the aprotinin coding sequence. Product was expressed as a fusion protein with beta-galactosidase. Cyanogen bromide cleavage of the fusion protein, followed by purification and renaturation is required to produce biologically active product.\nMarks, et al., in J. Biol. Chem. 261, pp. 7115-118 (1986), report the successful expression of bovine pancreatic trypsin inhibitor (BPTI) in E. coli. These workers were able to overcome the difficulties in producing disulfide-bonded proteins in E. coli, describing the production of correctly folded BPTI that was conformationally indistinguishable from native BPTI.\nOther publications dealing with the recombinant production of aprotinin include European Patent Application Nos. 238,993, 244,627, and 208,539.\nYeasts can offer clear advantages over bacteria in the production of heterologous proteins, which include their ability to secrete heterologous proteins into the culture medium. Secretion of proteins from cells is generally superior to production of proteins in the cytoplasm. Secreted products are obtained in a higher degree of initial purity; and further purification of the secreted products is made easier by the absence of cellular debris. In the case of sulfhydryl-rich proteins, there is another compelling reason for the development of eukaryotic hosts capable of secreting such proteins into the culture medium: their correct tertiary structure is produced and maintained via disulfide bonds. This is because the secretory pathway of the cell and the extracellular medium are oxidizing environments which can support disulfide bond formation [Smith, et al., Science, 229, 1219 (1985)]; whereas, in contrast, the cytoplasm is a reducing environment in which disulfide bonds cannot form. Upon cell breakage, too rapid formation of disulfide linkages can result in random disulfide bond formation. Consequently, production of sulfhydryl-rich proteins, such as APR, containing appropriately formed disulfide bonds, can be best achieved by transit through the secretory pathway.\nSecretion of aprotinin having the desired intramolecular disulfide bonds from E. coli is suggested in Japanese document No. 63230089 (assigned to Takeda Chemical Ind KK). The cited document contains no details as to the level of secretion or the purity of APR obtained, nor is there any suggestion that APR could be expressed in any other host system besides E. coli.\nPCT International Application No. PCT/DK88/00138, bearing International Publication No. WO 89/01968, describes the production of aprotinin and aprotinin homologues in S. cerevisiae under the control of the S. cerevisiae triosephosphate isomerase promoter. The aprotinin and aprotinin homologues are encoded by autonomously replicating plasmid-borne DNA. The reported yields of aprotinin and homologues thereof are quite low, ranging from 1-13 mg of aprotinin per liter of fermentation broth.\nIn view of the problems usually encountered with up-scaling the production of heterologous proteins in autonomous plasmid-based yeast systems, such as S. cerevisiae, and the low expression levels of aprotinin and aprotinin homologues achieved in reported work in S. cerevisiae, no motivation is provided by the art for one to further pursue the production of APR in S. cerevisiae. To overcome the major problems associated with the expression of recombinant gene products in S. cerevisiae (e.g., loss of selection for plasmid maintenance and problems concerning plasmid distribution, copy number and stability in fermentors operated at high cell density), a yeast expression system based on methylotrophic yeast, such as for example, Pichia pastoris, has been developed. A key feature of this unique system lies with the promoter employed to drive heterologous gene expression. This promoter, which is derived from a methanol-responsive gene of a methylotrophic yeast, is frequently highly expressed and tightly regulated (see, e.g., European Patent Application No. 85113737.2, published Jun. 4, 1986, under No. 0 183 071 and issued in the U.S. on Aug. 8, 1989, as U.S. Pat. No. 4,855,231). Another key feature of expression systems based on methylotrophic yeast is the ability of expression cassettes to stably integrate into the genome of the methylotrophic yeast host, thus significantly decreasing the chance of vector loss.\nAlthough the methylotrophic yeast P. pastoris has been used successfully for the production of various [Cregg et al., Bio/Technology 5, 479 (1987)], lysozyme and invertase [Digan et al., Developments in Industrial Microbiology 29, 59 (1988); Tschopp et al., Bio/Technology 5, 1305 (1987)], endeavors to produce other heterologous gene products in Pichia, especially by secretion, have given mixed results. At the present level of understanding of methylotrophic yeast expression systems, it is unpredictable whether a given gene can be expressed to an appreciable level in such yeast or whether the yeast host will tolerate the presence of the recombinant gene product in its cells. In addition, it is unpredictable whether desired or undesired proteolysis of the primary product will occur, and if the resulting proteolytic products are biologically active. Further, it is especially difficult to foresee if a particular protein will be secreted by the methylotrophic yeast host, and if it is, at what efficiency. Even for the non-methylotrophic yeast S. cerevisiae, which has been considerably more extensively studied than P. pastoris, the mechanism of protein secretion is not well defined and understood."}
{"text": "Typically such dispensing closures are designed such that the closures can be attached to a neck of a container, e.g. by screwing the dispensing closure onto the neck of a container, so that a fluid can be dispensed out of the container and through the opening or a spout of the dispensing closure.\nSuch dispensing closures are used for many different purposes and are frequently used for containers which contain fluid media, especially drinks or other method products, but such dispensing closures can be also used for containers storing other fluids or viscous products. Such a dispensing closure is for example disclosed in EP 1 503 942 B1."}
{"text": "In recent years many fabrication methods have been designed to form complicated shapes of resin compositions such as the hollow parts and tubes in various industrial applications. However, there are certain limitations to these existing methods.\nMany fabrication methods rely on adhesives for their sealing properties, but these are time-consuming and costly, and pose environmental concerns due to the use of volatile solvents. Ultrasonic welding or spin welding suffer from limitations on the shape and size of the objects bonded together, and occasionally show insufficient bonding strength. Vibration welding is often unattractive due to the inability to effectively control product appearance and flash, thereby limiting usage to certain applications.\nHence, the laser welding is increasingly attractive as a new method to better cope with these drawbacks. In laser welding, a laser light is irradiated through a transmitting resin material onto an absorptive resin material attaching to the resin material. The energy of the laser light accumulated on the contacting part of the absorptive resin material heats and melts the contacting part and the transmitting resin material is also heated and melted through heat transfer. The result of this operation is that the resin materials are easily and strongly joined together.\nAnother benefit to laser welding is that it increasingly offers freedom of choice in designing the shape of the joined articles because energy is applied in a noncontact fashion for the finishing product to be melted/bonded.\nSeveral important laser welding methods rely on Nd:YAG lasers (or known simply as YAG lasers) or diode lasers as the laser beam source, and these lasers emit light in the near infrared region. The diode laser techniques have become particularly advanced in recent years and diode lasers with higher output power can be obtained at lower cost.\nMany materials may benefit from welding techniques using such lasers. For example, polyethylene resin, polypropylene resin, polystyrene resin, polycarbonate resin, acrylic resin and nylon resin have been demonstrated as effective candidates for laser welding. Thermoplastic resin compositions useful in laser welding are described, for example, in Japanese Published (Koukoku) Patent No. 62-49850 and Japanese Published (Koukoku) Patent No. 5(93)-42336. Other resin compositions associated with the laser welding are described in U.S. Pat. No. 5,893,959 in which carbon black or nigrosine is used as a colorant for thermoplastic resin.\nIf there are many efforts directed to the laser welding of nylon resins. In conventional laser welding, laser beams penetrate through a laser transmitting article positioned close to a laser beam source, and are largely absorbed in the laser absorbing article disposed in contact with the laser transmitting article. This causes the junction portion to be melted and jointed together. However, non-colored resins have been used as the transmitting resin material. The use of such materials limits their applicability for articles of various colors demanded in the automotive industry and electric/electronic industries. Of particular interest, the use of black material in these applications is not satisfactorily popularized at this time using conventional laser welding operations. Additionally, there are some suggestions that black pigment can be diluted and utilized in part of the transmitting resin or even using materials in a thinner shape to facilitate transmission. However such approaches cannot ensure the satisfactory appearance of the resulting part and do not allow much flexibility in designing parts. There are still other examples suggesting the addition of carbon black to the absorptive resin as an approach. However the details of such an approach are not yet fully understood or functional.\nThe present invention provides a thermoplastic resin composition capable of offering molded articles which appear in black, are transparent to a laser beam at wavelengths in the infrared region. As another feature, it provides a substantially homogenous visual black impression in combination with opaque articles that appears in black and absorbs the laser beam largely by containing black dyes, welded together by the laser beam. These materials offer advantages in excellent and balanced heat-resistance and mechanical properties as required in automotive parts, electric/electronic components, mechanical components, and many other applications. These and other objects, features and advantages of the invention will become better understood upon having reference to the following description of the invention."}
{"text": "This disclosure relates to communications systems and, more particularly, to an improved spontaneous messaging service.\nSpontaneous messaging is a communications service feature that is invoked when a communications service provider\"\"s network receives some indication that a call cannot be completed, due to a ring-no-answer or a busy-line condition at a called party\"\"s telephone set. The network or a component thereof, spontaneously invites the caller to record a message for the called party. Such a message is hereinafter referred to as a xe2x80x9cspontaneous messagexe2x80x9d. The recorded message is initially stored and subsequently delivered to the called party upon his or her availability. Specifically when the caller opts to record a spontaneous message for the called party, the network or a component thereof, attempts to deliver such spontaneous messages every thirty (30) minutes for example, for up to a certain number of hours. Spontaneous messaging service is presently available for credit card, calling card and conventional calls from communications carriers such as ATandT.\nOver the last few years, however an increasing number of communications services users have evinced a preference for using pre-paid calling cards, as opposed to calling cards to complete telephone calls. Users\"\" preference for pre-paid calling cards over conventional calling cards can be attributed to a number of factors that include elimination of bill payment, certainty about the cost per minute for the call, cheaper rates than calling card calls and ease of use, to name a few. Such ease-of-use is evidenced by the fact that a pre-paid calling card user who wishes to complete a calling card call needs simply to dial a toll-free telephone number, and then key in the prepaid calling card number (and optional password) printed on the prepaid calling card) followed by the destination telephone number. The call is then completed by the service provider that issued the pre-paid calling card after the appropriate checks for the availability and sufficiency of a balance on the account associated with the pre-paid calling card number. At predetermined time intervals during the call, the service provider debits from the account associated with the pre-paid calling card the appropriate amount of calling time used by the caller. Because the pre-paid card billing system deducts a fixed rate only for completed calls or other services rendered, pre-paid calling card users do not pay for services when the call is not completed due to ring-no-answer or busy-line conditions.\nIn spite of the popularity of pre-paid calling cards, the spontaneous messaging service feature is still unable to allow a pre-paid card caller to record a message for an unavailable called party. This is because the spontaneous messaging service has no mechanism to deliver such service for a fee. Specifically, the spontaneous messaging service has to identify an unequivocal billing number before a caller is allowed to record a spontaneous message. Such billing number is used to charge the caller when the recorded spontaneous message is delivered to the called party. In the case of pre-paid calling cards, a series of factors operate to make this conventional billing technique unavailing. One such factor is the uncertainty with respect to the availability and sufficiency of a balance on the account associated with the pre-paid card number to pay for message delivery when it occurs. For example, the account associated with the pre-paid card number may have been depleted by the time the message is delivered to the called party. Equally problematic is the lack of a mechanism to credit pre-paid calling card accounts for undelivered messages. Specifically, if the pre-paid card caller is charged for delivery of a message that cannot be delivered within the pre-determined delivery time window, the pre-paid caller would be charged for service that has not been rendered. Thus, a problem of the prior art is the inability of pre-paid calling card callers to receive spontaneous messaging service.\nFrom a general perspective, prepaid calling card users are unable to receive any communication service that requires a refund mechanism to credit such users for any service charge when the communication service cannot be delivered.\nAt a yet higher level of generality, debit card users are unable to receive any service that may require an electronic refund of the fees charged when such service cannot be delivered within a predicted period of time.\nThe present disclosure is directed to a method and a system that allow an account associated with a debit card to be debited to an escrow account when the user of such debit card elects to receive a service that is anticipated to be delivered or provided within a predetermined period of time. For example, an account associated with a pre-paid calling card number will be debited when a pre-paid calling card caller records a spontaneous message for a called party. The debited amount in the escrow account is either permanently deducted from the prepaid calling card account when the message is delivered, or partially or fully credited to the prepaid calling card account when the recorded message is deemed undeliverable.\nIn an exemplary embodiment of the principles disclosed herein, a pre-paid calling card caller is prompted to record a spontaneous message when the caller places a call to a called party who is unavailable due to a ring-no-answer or busy-line condition at the called party\"\"s telephone set. If the pre-paid calling card caller indicates a desire to record a spontaneous message for the unavailable called party, the service provider\"\"s network performs certain functions before such recording. One such function is to ascertain whether the remaining balance of the account associated with the pre-paid calling card number is sufficient to defray the anticipated cost of the spontaneous messaging service. Another one of such functions is to inform the caller that the cost of the service would be fully or partially refunded in the event that the message cannot be delivered to the called party within a predetermined time period. Yet another function performed by the network of the service provider is to create an electronic escrow account for the card number. Finally, the service provider\"\"s network transfers the message delivery service fee from the account of the pre-paid calling card number to the newly created escrow account.\nIf the spontaneous message previously recorded by the caller is delivered to the called party, the amount of money in escrow is then transferred to the card issuer. When the messaging service is unable to deliver the spontaneous message to the called party before the expiration of the predetermined time period for such delivery, then the amount of money in escrow is returned to the account for the pre-paid calling card (less any agreed-upon usage fee, if any)."}
{"text": "Different functional layers in different organic semiconductor devices request a variety of special characteristics.\nFor instance organic thin-film transistors (OTFTs) need high mobility materials in their active channel. Transparent circuits, such as transparent OTFTs require that the high mobility organic material also comprises a wide electronic band gap; the electric injection of holes and/or electrons must be still provided.\nOLEDs require transparent transport layers, with high conductivity. The transparency is necessary in those opto-electric devices to avoid non desired absorption of the light. These so called “window” materials can be used as transport layers, exciton or charge blocking layers. The thickness of the layers made with the window materials is used to adjust the micro cavity of the OLEDs in such a way that the outcoupled emission of the OLED is a maximum. The non-optically active layers of all kinds of semiconductor devices can be exchanged for window materials in order to fabricate fully transparent components and circuits (e.g US20060033115). The functionality and nomenclature of the layers are typical as used in the field. Further explanation can be found in US2006244370.\nElectronic devices also need high stability towards temperature, meaning that the intrinsic properties of the amorphous organic semiconducting materials, such as triphenyl amine derivatives, or phenantronine derivatives, must include a high glass transition temperature (Tg) and high temperature stability in the device.\nThe performance characteristics of (opto)electronic multilayered components are determined by the ability of the layers to transport the charge carriers, amongst others. In the case of light-emitting diodes, the ohmic losses in the charge transport layers during operation are associated with their conductivity. The conductivity directly influences the operating voltage required and also determines the thermal load of the component. Furthermore, depending on the charge carrier concentration in the organic layers, a bending of the band in the vicinity of a metal contact results which simplifies the injection of charge carriers and can therefore reduce the contact resistance.\nBy electrically doping hole transport layers with a suitable acceptor material (p-doping) or electron transport layers with a donor material (n-doping), respectively, the density of charge carriers in organic solids (and therefore the conductivity) can be increased substantially. Additionally, analogous to the experience with inorganic semiconductors, applications can be anticipated which are precisely based on the use of p- and n-doped layers in a component and otherwise would be not conceivable. The use of doped charge-carrier transport layers (p-doping of the hole transport layer by admixture of acceptor-like molecules, n-doping of the electron transport layer by admixture of donor-like molecules) in organic light-emitting diodes is described in US2008203406 and U.S. Pat. No. 5,093,698.\nUS2008227979 discloses in detail the doping of organic transport materials, also called matrix, with inorganic and with organic dopants. Basically, an effective electronic transfer occurs from the dopant to the matrix increasing the Fermi level of the matrix. For an efficient transfer in a p-doping case, the LUMO energy level of the dopant must be more negative than the HOMO energy level of the matrix or at least slightly more positive, not more than 0.5 eV, to the HOMO energy level of the matrix. For the n-doping case, the HOMO energy level of the dopant must be more positive than the LUMO energy level of the matrix or at least slightly more negative, not lower than 0.5 eV, to the LUMO energy level of the matrix. It is furthermore desired that the energy level difference for energy transfer from dopant to matrix is smaller than +0.3 eV.\nTypical examples of doped hole transport materials are: copperphthalocyanine (CuPc), which HOMO level is approximately −5.2 eV, doped with tetrafluoro-tetracyanoquinonedimethane (F4TCNQ), which LUMO level is about −5.2 eV; zincphthalocyanine (ZnPc) (HOMO=−5.2 eV) doped with F4TCNQ; a-NPD (N,N′-Bis(naphthalen-1-yl)-N,N-bis(phenyl)-benzidine) doped with F4TCNQ.\nTypical examples of doped electron transport materials are: fullerene C60 doped with acridine orange base (AOB); perylene-3,4,9,10-tetracarboxylic-3,4,9,10-dianhydride (PTCDA) doped with leuco crystal violet; 2,9-di(phenanthren-9-yl)-4,7-diphenyl-1,10-phenanthroline doped with tetrakis(1,3,4,6,7,8-hexahydro-2H-pyrimido[1,2-a]pyrimidinato)ditungsten (II) (W(hpp)4); naphthalene tetracarboxylic acid di-anhydride (NTCDA) doped with 3,6-bis-(dimethyl amino)-acridine; NTCDA doped with bis(ethylene-dithio)tetrathiafulvalene (BEDT-TTF).\nThere is a technical challenge to provide electron transport materials (ETM) and emitter host (EMH) materials that have a sufficiently low laying LUMO level so that they can be doped, and still have a high enough laying LUMO level which can efficiently transfer charge to emitter host (in case of an ETM) and transfer energy to the emitter dopant (in case of EMH). The limitation for high laying LUMO level of the ETL is given by the dopability, since the n-dopants with very high HOMO tend to be unstable; also the injection is difficult for very high LUMO of the ETL."}
{"text": "1. Field of the Invention\nThis invention relates to the field of specific binding assays, particularly immunoassays for determining substances of clinical interest. The development of specific binding assay techniques has provided extremely useful analytical methods for determining various organic substances of diagnostic, medical, environmental and industrial importance which appear in liquid mediums at very low concentrations. Specific binding assays are based on the specific interaction between a ligand, i.e., a bindable analyte under determination, and a binding partner therefor, i.e., receptor. The presence of the receptor can be used to effect a mechanical separation of bound and unbound labeled analyte or can affect the label in such a way as to modulate the detectable signal. The former situation is normally referred to as heterogeneous and the latter as homogeneous, in that the latter technique avoids a separation step. Where one of the ligand and its binding partner is a hapten or antigen and the other is a corresponding antibody, the assay is known as an immunoassay. See, generally, Odell and Daughaday (Eds.), Principles of Competitive Protein-Binding Assays, J. B. Lippincott Co., Philadelphia (1971).\nIn conventional label conjugate specific binding assay techniques, a sample of the liquid medium to be assayed is combined with various reagent compositions. Such compositions include a label conjugate comprising a binding component incorporated with a label. The binding component in the conjugate participates with other constituents, if any, of the reagent composition and the ligand in the medium under assay to form a binding reaction system producing two species or forms of the conjugate, e.g., a bound-species (conjugate complex) and a free-species. In the bound-species, the binding component, e.g., a hapten or antigen, of the conjugate is bound by a corresponding binding partner, e.g., an antibody, whereas in the free-species, the binding component is not so bound. The relative amount or proportion of the conjugate that results in the bound-species compared to the free-species is a function of the presence (or amount) of the ligand to be detected in the test sample.\n2. Brief Description of the Prior Art\nMany properties of natural cell membranes can be duplicated in simple lipid bilayer systems, referred to as liposomes. One of these properties is lysis. Lysis can be achieved by conventional chemical agents, such as detergents, or by immunological reactions. When a vesicular, e.g., liposome, membrane contains an externally accessible antigen it will react with corresponding antibody, causing agglutination. When the antigen-sensitized liposome reacts with corresponding antibody in the presence of complement the membrane is irreversibly damaged and can no longer function as the intact selective permeability barrier. This is immunolysis. The extent of immunolysis has been monitored by using antigen-sensitized liposomes containing entrapped marker molecules which are released upon immunolysis.\nThe lysis of liposomes has been studied using a wide variety of marker systems. One type of marker system, which uses no marker-reactive reagents external to the liposome, has been disclosed in which the liposome encloses an electron paramagnetic resonance spin marker, e.g., stable free radical, such as tempocholine. Tempocholine is water soluble but membrane impermeable, therefore it can be entrapped within liposomes or vesicles and it does not leak across the lipid bilayers. Tempocholine enclosed within liposomes produces a characteristic broad paramagnetic resonance signal of small amplitude. This is because there is a high concentration of spin molecules which exchange signals due to their close proximity to one another. When the liposome membranes are ruptured, such as by activation of complement, the tempocholine molecules are released and diluted in the external medium. This results in a readily detectable, qualitative and quantitative, alteration in the paramagnetic resonance spectrum. Humphries, G. M. and McConnell, H. M. Proc. Nat. Acad. Sci. USA, 72: 2483-2487 (1975). See, also, Wei, et al., J. Immunol. Methods, 9: 165-170 (1975); Chan, et al., J. Immunol. Methods, 21: 185-195 (1978); and Hsia, et al., Ann. N.Y. Acad. Sci., 308: 139-148 (1978). Another marker system which uses no marker-reactive reagents external to the liposome employs a composition containing a fluor, such as 1-aminonaphthalene-3,6,8-trisulfonate, and a quencher, such as a,a'-dipyridinium p-xylene dibromide, entrapped within the liposome. Fluor and quencher escape upon immunolysis of the liposome and their subsequent dilution in the external volume abolishes the quenching, resulting in a high fluorescent signal. See Smolarsky, et al., J. Immunol. Methods, 15: 255-265 (1977) and Geiger, et al., J. Immunol. Methods, 17: 7-19 (1977).\nAnother type of marker system employs an electrode in the reaction environment external to the liposome. One example of this uses a potassium loaded liposome. Potassium ion escapes upon liposome immunolysis and reacts with an ion-selective electrode in the external environment. See Katsu, et al., Chem. Pharm. Bull, 30: 1504-1507 (1982). Tetrapentylammonium ions have also been enclosed in liposomes and detected upon immunolysis using an ion-selective electrode. See Shiba, et al., Anal. Chem., 52: 1610 (1980) and Chem. Lett. 155 (1980). Also, sheep erythrocyte ghosts (membranes) have been loaded with trimethylphenyl ammonium ion and lysis has been detected with an ion-selective electrode. D'Orazio, et al., Anal. Chem., 49: 2083 (1977) and D'Orazio, et al., Anal. Chem. Acta. 25: 109 (1979).\nOne of the first types of marker systems to be reported employs substrate entrapped within the liposome or erythrocyte membrane vesicle. The substrate escapes upon immunolysis of the liposome and reacts with an enzyme-containing composition in the external volume to produce a detectable response. The early examples of this use glucose entrapped within the liposome. The glucose escapes upon immunolysis and its release from the liposomes is measured by the increase in absorbance at 340 nanometers which occurs upon reduction of NADP.sup.+ in the presence of hexokinase, glucose-6-phosphate dehydrogenase and the requisite cofactors. See Hixby, et al., Proc. Nat. Acad. Sci., 64: 290-295 (1969); Kinsky, et al., Biochemistry, 8: 4149-4158 (1969); Kinsky, et al., Biochemistry, 9: 1048 (1970). In a more recent example of this type of system, a fluorogenic substrate (umbelliferone phosphate) or a chromogenic substrate (p-nitrophenyl phosphate) is enclosed in the liposome. The substrate escapes upon immunolysis of the liposome and reacts with an enzyme (alkaline phosphatase) in the external volume. Free substrate is produced, resulting in an increased signal. See Six, et al., Biochemistry, 13: 4050 (1974); Uemura, et al., J. Biochem, 87: 1221 (1980); and Uemura, et al., J. Immunol. Methods, 53: 221-232 (1982).\nLysis, including immunolysis, of liposomes having internally entrapped enzymes as markers has also been reported. The enzyme reacts with substrate in the external reaction medium resulting in a detectable response. For example, experiments have been performed using trapped enzyme markers including hexokinase, glucose-6-phosphate dehydrogenase and B-galactosidase. See Kataoka, et al., Biochem. Biophys. Acta, 298: 158-179 (1973). Also, horseradish peroxidase can be entrapped within the liposome. The peroxidase escapes upon immunolysis and catalyzes the following reaction: ##STR1## Oxygen is consumed by the oxidation of NADH and resultant production of water. The depletion of oxygen is detected by an oxygen electrode. See Haga, Biochem. Biophys. Res. Comm., 95: 187-192 (1980).\nSeveral references report the enhancement of enzymic activity upon lysis of enzyme-containing liposomes. For example, Solomon, et al., Biochem. Biophys. Acta, 455: 332-342 (1976) report this observation upon lysis of glucose oxidase-containing liposomes by either sonication or exposure to detergent. The enzymic activity of the entrapped glucose oxidase served as a measure for the permeability of the bilayer membrane of the liposomes to glucose in a non-separation assay. The oxidation of glucose was followed in the same reaction mixture before and after detergent lysis of the enzyme-containing liposome by oxygen uptake using an O.sub.2 electrode. Observed enzymic activity was as high as 4.60 times greater after lysis than it was before lysis.\nTokunaga, et al., FEBS Letters, 106: 85-88 (1979) report an in vitro non-separation method for determining permeability of liposomes containing alkaline phosphatase, a-glucosidase and a-galactosidase. Enzymic activity was determined using substrate and a coupled enzyme reaction in the same reaction mixture before and after detergent lysis. Using alkaline phosphatase, observed enzymic activity was as high as 17.5 times greater after lysis than it was before lysis. These authors state that if the K.sub.m of the particular enzyme/substrate pair were the same both inside and outside of the liposome, the intravesicular enzyme activity could be related to the apparent rate of substrate permeability.\nMagee, et al., J. Cell Biology, 63: 492-504 (1974) have used horseradish peroxidase-containing liposomes. The enzymic activity of horseradish peroxidase was determined spectrophotometrically in the same reaction mixture before and after detergent lysis. The ratios of the observed enzymic activity after lysis to the observed enzymic activity before lysis were as high as 8.3. They note that liposomes have been reported as useful models for membranes in permeability studies owing to their sensitivity to polyene antibiotics and susceptibility to immune lysis.\nSeveral different homogeneous label conjugate specific binding assay systems are known in the art, one of which is disclosed by Ullman, et al., in U.S. Pat. No. 4,193,983. In this assay system, a label and a ligand or ligand analog are non-covalently bound to the external surface of a colloidal particle, such as a liposome, which is capable of maintaining its integrity in an aqueous environment. The discrete colloidal particle serves as a hub or nucleus which retains the ligand or its analog and the label in a substantially fixed average spatial relationship. By having the label and ligand in relatively close proximity on the surface of the particle, the proximity of the label and the receptor bound to the ligand adjacent the label can be used in accordance with the prior art techniques to modulate the signal from the label. Nowhere is immunolysis suggested. Indeed, the clear teaching of this reference is that lysis of the vesicle would adversely affect the spatial relationship (proximity) of label and ligand which is required for this assay procedure.\nSpecific binding assay systems have been proposed, using a multilayered lipid membrane vesicle which has been prepared or treated to have surface-bound ligand or ligand analog and a marker or reagent substance enclosed within the vesicle. The remaining reagents for the assay include: (1) a binding partner, e.g., antibody, for the ligand; and (2) complement to effect lysis of the vesicle upon binding of the binding partner to surface-bound ligand. Generally, see McConnell, U.S. Pat. No. 3,850,578 and McConnell, et al., U.S. Pat. No. 3,887,698 and Gregoriadis, et al., Liposomes in Biological Systems, John Wiley & Sons, N.Y. (1980), especially Chapter 12 entitled \"Liposomes as Diagnostic Tools\".\nMore particularly, immunoassay systems have been disclosed in which the use of enzyme-encapsulating liposomes is suggested. Hsia, et al., U.S. Pat. No. 4,235,792 describes a competitive homogeneous immunoassay method which employs immunolysis of an antigen-sensitized liposome containing a marker. Enzymes are among the markers disclosed (col. 6, lines 24-28).\nCole, U.S. Pat. No. 4,342,826 discloses a specific binding assay using antigen-sensitized, enzyme-containing liposomes. These liposomes are immunospecifically caused to release enzyme upon binding of corresponding antibody and fixing of active complement. Upon enzyme release, the presence or absence of enzymatic activity is detected. Cole emphasizes the advantage of providing a homogeneous system in which enzymic activity is substantially greater upon lysis, e.g., a \"signal:noise\" ratio of at least 5-10 and preferably above 60.\nEach of the above approaches to vesicular marker systems has provided an advance of one sort or another in sequestering marker from the reaction medium and, thus, minimizing the generation of signal (latency) prior to immunolysis. That is, the signal observed from intact liposomes is considerably less than that from lysed liposomes. As demonstrated by the references, this end has been widely recognized as a major consideration in the improvement of liposome immunoassays. Further, the combined teaching of the literature in this area has been that the advantages to be achieved are enhanced when complete sequestration prior to lysis is combined with the production of a high intensity signal upon immunolysis."}
{"text": "1. Field of the Invention\nThe present invention relates to a flash memory, and more particularly, to a self-test method of a flash memory device.\n2. Description of the Prior Art\nA flash memory can store data through operations of electronically erasing and writing/programing. The flash memory is widely applied to electronic products such as memory cards, solid-state drives (SSDs) and portable multimedia players. Since the flash memory is a non-volatile memory, it does not need additional power to keep the information stored therein. Further, the flash memory is capable of providing fast data read and better shock-resistive ability. Hence, based on at least the above characteristics, the reason why the flash memory is so popular can be explained.\nFlash memories can be divided into NOR flash memories and NAND flash memories. The NAND flash memory has shorter erase and program time, and each memory cell thereof requires a smaller chip area. Hence, compared with the NOR flash memory, the NAND flash memory has higher storage density and lower cost per storage unit. In general, the flash memory utilizes memory cell arrays to store data bits. The memory cell is implemented by a floating-gate transistor, and each memory cell may set the required threshold voltage of turning on the floating-gate transistor through appropriately controlling the number of charges on the floating gate of the floating-gate transistor, to store the information of a single bit or multiple bits. In this way, when one or multiple predetermined control gate voltages are applied to the control gate of the floating gate transistor, the conduction state of the floating-gate transistor will indicate one or multiple binary digits stored in the gloating-gate transistor.\nHowever, some factors might affect/disturb the number of charges originally stored in the flash memory cell. For example, the interference existing in the flash memory may come from the write/program disturbance, read disturbance and/or retention disturbance. Taking a NAND flash memory composed of memory cells each having more than one bit stored therein for example, a physical memory page includes multiple logical memory pages, and each of the logic memory pages utilizes one or multiple control gate voltages to perform a read operation. For example, regarding a flash memory cell used for storing 3-bit information, the flash memory cell will have one of 8 states (i.e., charge levels) respectively corresponding to different numbers of charges (i.e., different threshold voltages). However, due to the program/erase (P/E) count and/or data retention time, the threshold voltage distribution of the flash memory cell will change. Hence, reading information stored in the memory cell according to the original control gate voltage setting may not correctly obtain the stored information due to the change of the threshold voltage distribution.\nThe aforementioned change or shift of the threshold voltage distribution usually results in data read errors, and this problem can be improved by utilizing error correction code (ECC) and soft decoding. However, due to the reduction on the current manufacture process of the flash memories, it becomes easier to encounter errors during the procedure of soft recording preformed on data read from the flash memories. Hence, how to provide a low-cost and efficient testing method becomes an important issue in the pertinent field."}
{"text": "1. Field of the Invention\nThe invention relates to a method of blocking aberrant Ras signaling in a mammal while avoiding excessive cell toxicity.\n2. Background Information\nThe mevalonic acid (MVA) pathway is responsible for the biosynthesis of cholesterol and isoprenoid intermediates such as geranylgeranylpyrophosphate (GGPP) and farnesylpyrophosphate (FPP). Two sites in the MVA pathway have been cited to be of particular importance: the synthesis of MVA by 3-hydroxy-3-methylglutaryl-coenzyme A reductase (HMG-CoA reductase), an early step thought to be the major point of regulation, and the so-called \"branch-point\" of FPP metabolism (Brown and Goldstein 1980; Sabine, 1983; reviewed in Grunler et al., 1994). FPP is the last common intermediate in the pathway and is the substrate for a number of different enzymes that catalyze committed steps in branching pathways leading to the biosynthesis of cholesterol, ubiquinone, dolichol, as well as isoprenylated proteins and hemes. GGPP synthase, one of the branch-point enzymes, catalyzes the condensation of FPP and isopentenyl pyrophosphate to form GGPP. GGPP and FPP are utilized by geranylgeranyl-transferases (GGTases) I and II, and farnesyltransferase (FTase), respectively, for posttranslational isoprenylation of proteins on carboxyl terminal cysteine residues (reviewed in Maltese, 1990; Casey, 1992; Grunler et al., 1994). FTase and GGTase I prenylate proteins with carboxyl termini that end with a CAAX box where C=cysteine, A=aliphatic, and X=any amino acid. FTase prefers X as a serine or methionine whereas GGTase I prefers X as a leucine or isoleucine. GGTase II prenylates proteins that end in XXCC and XCX where X is any amino acid. For several proteins, isoprenylation is essential for proper intracellular localization and biological function (Holtz et al., 1989; Fukada et al. 1990; Der and Cox 1991; Hori et al. 1991; Inglese et al. 1992). In contrast to FPP, GGPP is currently known to be utilized only for protein prenylation.\nGeranylgeranylated proteins and farnesylated proteins appear to comprise distinct but overlapping, sets of proteins, with the former being greater in number than the latter (Farnsworth et al., 1990; Epstein et al., 1990). Many of these proteins have been shown to play essential roles in signal transduction pathways and some have been implicated in malignant transformation. For example, the geranylgeranylated low molecular weight (20-28 KDa) guanine nucleotide-binding proteins Rho and Rac have recently been shown to be critical players in regulating not only the organization of the actin cytoskeleton (Nobes and Hall 1995) but also the progression of the cell cycle through Gl (Olson et al. 1995). In addition Ras, another family of guanine nucleotide-binding proteins, control normal cell growth (Mulcahy et al., 1985) and differentiation (Bar-Sagi and Feramisco, 1985) and, when mutated, can produce malignant transformation (Reddy et al., 1982). The Ras family of proteins serve as transducers of extracellular signals from receptor tysosine kinases to the nucleus (McCormick 1993). Their stimulation by these receptors results in the activation of several growth-related pathways including a cascade of mitogen-activated protein (MAP) kinases such as Raf, MEK and ERK (McCormick 1993), the latter of which can translocate to the nucleus and regulate the activity of some transcription factors. In some human cancers, Ras is GTP-locked and constitutively activates the MAPK cascade. Such cancers include, but are not limited to, colorectal, pancreatic and lung carcinomas, and melanoma.\nThe ability of Ras to cause cancer requires its attachment to the plasma membrane which is mediated by prenylation (Der and Cox 1991; Kato et al., 1992). The prenylation of Ras in vitro exhibits a preference for farnesylation (James et al, 1995). Furthermore although the prenylation of K.sub.B -Ras in vivo is, at present, uncertain (Casey et al., 1989; Lerner et al., 1995b), H-Ras has been shown to require farnesylation for its cancer-causing activity (Hancock et al. 1989; Seabra et al., 1991; Lerner et al. 1995a). Moreover, in nude mice, inhibitors of FTase are effective at suppressing the growth of tumor cells possessing oncogenic H- or K-Ras (Sun et al., 1995).\nLovastatin is a potent competitive inhibitor of HMG-CoA reductase (Alberts 1988) and is used clinically as a cholesterol lowering agent. Following oral administration, the inactive lactone form is hydrolyzed to the .beta. hydroxy acid form, which inhibits HMG-CoA reductase. HMG-CoA reductase catalyzes the reduction of HMG-CoA to mevalonate, an early rate-limiting step in the biosynthesis of cholesterol.\nLovastatin significantly reduces not only the biosynthesis of the end-product cholesterol for which it is used clinically (Illingworth and Bacon 1989), but also depletes the intracellular pools of GGPP and FPP, resulting in the inhibition of protein geranylgeranylation and protein farnesylation (Leonard et al., 1990). It has been recently discovered that lovastatin disrupts early signaling events such as tyrosine phosphorylation levels of the PDGF receptor and its association with PI-3-kinase (McGuire et al., 1993). Moreover, lovastatin arrests cultured cells predominantly in the G.sub.1 phase of the cell cycle and produces a characteristic rounded morphology (Quesney-Hunecus et al., 1979; Sinensky and Logel, 1985; Fenton et al., 1992). Lovastatin has also been found to inhibit tumor growth of cells expressing oncogenic H-Ras in nude mice, but at doses that blocked tumor growth, the animals died (Sebti et al., 1991). It is not known whether inhibition of protein farnesylation and/or protein geranylgeranylation, or the reduction in levels of some other end-product of the MVA pathway, are responsible for lovastatin's toxicity."}
{"text": "Outdoor mailboxes typically are positioned close to roads. Therefore, such a mailbox is susceptible to being hit by a motor vehicle, or by other objects such as snow plowed by a snowplow, or even by vandals. The mailbox post is then vulnerable to being broken.\nIt is known to configure a mailbox post so that it has a flexible joint, thereby allowing the mailbox post to bend rather than break. See Dunn (U.S. Pat. No. 6,223,982). Dunn discloses installing a non-circular spring wire between the upper and lower parts of the mailbox post, the spring being positioned vertically, so that the mailbox post will return to a vertical position after being hit.\nHowever, there are several drawbacks to Dunn's configuration. First, Dunn does not disclose that the spring unit is distinct from the post, and thus if any part of Dunn's post becomes damaged then the entire thing (including the joint) either becomes useless or must be repaired. Likewise, the invention of Dunn replaces an existing mailbox post, instead of being adaptable to an existing mailbox post.\nAlso, Dunn requires a non-circular elastic element in order for the mailbox post to properly align with a horizontal axis, but such a spring is not as strong and durable as a circular spring. This non-circular elastic element also makes assembly and disassembly difficult, while putting great stress on the parts that secure the non-circular spring to the rest of the mailbox post. Moreover, the tension of Dunn's elastic piece is not adjustable."}
{"text": "1. Field of the Invention\nThe present invention relates to an intake and exhaust gas claning apparatus, and more particularly to an exhaust gas cleaner for automobiles, motorcycles and factories.\n2. Description of the Prior Art\nAir pollution is becoming a serious problem worldwide. Engine exhaust of automobiles and motorcycles, and factory discharges are the major sources of pollutants. Engine exhaust pollutants of a motor vehicle are generated mostly by the following causes:\n1. The low temperature in the combustion chamber when the engine is starting from cold. PA1 2. Driving uphill, bumpy road conditions or heavy loads. PA1 3. Poor quality of gasoline or oils. PA1 4. Clogged air filter. PA1 5. Intake of poor quality air from a tunnel or congested road. PA1 6. Engine abrasion or too many pollutants in the combustion chamber. PA1 7. Back pressure caused by a smoke filter or catalytic converter. PA1 1. High volume and high temperature. PA1 2. High speed and impulsive vibration. PA1 3. Viscosity after dispersion. PA1 4. High vibration frequency and noise level. PA1 5. Contains fine particles and becomes greasy. PA1 6. Contains water vapor. PA1 7. Contains gases of carbon dioxide, sulfer dioxide, carbon monoxide, etc. PA1 8. Contains polynuclear aromatic hydrocarbons (PAH) and other cancer inducing or cell mutation gases. PA1 9. Indirectly forms peroxyacety/nitrate (PAN) and other photochemicals which are environmentally harmful. PA1 1. Wet type dust collector: tends to be blown away by vapor, consumes much water, tends to form erosive gas, not economical or pratical. PA1 2. Honeycomb filter gauze formed of ceramic or metal: easily clogged by greasy smoke, incur back pressure, difficult operation, easily damaged, high cost, time consuming, short effective life, requires considerable backup materials. PA1 3. Whirlwind type or static collector: large size, consumes much electric power, heavy load, not suitable for motor vehicles. PA1 4. Additives in gasoline or diesel oil: while there are claims that the additives can improve combustion, not much evidence or proof is available for their effectiveness in eliminating toxic gas, and not conclusive as to whether or not the additives will be harmful to the engine. PA1 5. Combustion smoke filter: provides better results, but high cost, consumes energy and oxygen, and produces more carbon dioxide gas, has explosion risk, limited effective life, not recyclable.\nEngine exhaust gas includes smoke which contains solid particles and liquid drops. They all create severe air pollution problems and threaten the health of human beings and animals. Some of the exhaust gases discharge into the atmosphere and form \"acid rain\" which further creates a water pollution problem. The carbon dioxide in engine exhaust gas is also one of the factors causing the global greenhouse effect. Thus, to minimize engine exhaust pollution is becoming a worldwide focus.\nEngine exhaust gases have the following characteristics:\nMany types of engine exhaust controlling or eliminating means have been developed, but they have disadvantages, as indicated below:\nThe above disadvantages have limited conventional exhaust controlling devices from being widely used. Thus, there has been a need for providing a more cost effective engine exhaust controlling means."}
{"text": "1. Field of the Invention\nThis invention relates to a liquid composition for metachromatic members, and a metachromatic member set making use of the composition. More particularly, it relates to a liquid composition for metachromatic members which is used to form desired images on metachromatic members, and a metachromatic member set making use of the composition.\n2. Related Background Art\nConventionally, as an attempt to form images on metachromatic members, a doll- or animal-shaped toy set is disclosed which comprises a doll- or animal-shaped toy and a water-metachromatic member comprising a support and provided thereon a porous layer formed by fixing therein a low refractive index pigment in a disperse state together with a binder resin (see, e.g., Japanese Patent Application Laid-open No. 11-216272).\nIn the above doll- or animal-shaped toy set, any desired portion of a doll- or animal-shaped toy in a liquid-absorbed state is pressed against the surface of the water-metachromatic member to make the porous layer permeated with water and turn transparent or semitransparent so as to make an image visible. The toy has the repetition utility that, upon evaporation of the liquid having permeated the porous layer, the porous layer returns to the original state and the image comes invisible.\nIt is also disclosed therein that water is preferably applied as a medium which makes the toy liquid-absorbed and that, in order to control drying speed to elongate image visual time, a water-soluble organic solvent such as propylene glycol may be mixed in a very small quantity.\nHowever, the image formed in the water-metachromatic member constituted as described above is one which turns invisible upon evaporation of the water and, even when the water-soluble organic solvent such as propylene glycol is added, the image visual time is a little elongated at best.\nHence, such a means is unsuitable for uses in which images are permanently formed and preserved, e.g., uses in which impressions of hands, feet and the like of infants are taken and preserved. Accordingly, a liquid composition suited for the preservation of images has eagerly been awaited."}
{"text": "The present invention relates to an improved drum assembly for washing, drying machines, and the like, as well as a method for making said drum assembly.\nThe drum assemblies of the type concerned by the present invention conventionally comprise a body of substantially drum configuration. This drum body is formed by a perforated metal band constituting the cylindric wall of the drum body, by a rear closure disc for closing the drum body, which supports the cross member supporting the shaft of the drum, and a front ring element defining the front loading mouth therethrough the cloth articles to be processed are loaded into the drum. In the case of a top loading, on the other hand, the front ring element is replaced by a closing cover, whereas the side band is provided with an opening for introducing the cloth articles into the drum.\nIn order to provide an efficient tumbling of the cloth articles to be washed or dried, on the inside surface of the perforated metal band forming the cylindric wall of the drum are provided a plurality of projections, having a substantially wedge-like cross section, which operate as xe2x80x9centrainmentxe2x80x9d elements. Actually, owing to the provision of these projections, the rotary movement applied to the drum assembly will cause the linen loaded therein to be upwardly entrained and fall by gravity, thereby aiding the washing operation.\nThe provision of the above mentioned projections, which will be hereinafter called xe2x80x9cwedge elementsxe2x80x9d, causes, however, problems or difficulties in properly coupling the metal band and ring element or front cover. Because of the provision of the mentioned wedge elements, in fact, the edge or rim of the mentioned metal band will have an uneven contour, which, on the other hand, must be perfectly fitted or mated with the contour of the ring element or front cover of the drum.\nPrior solutions for providing such a connection include the use of spot welding operations between the respective walls of the ring or front cover element and the drum wedge elements, formed on the corresponding surface engagement portions. Such a spot welding method, however, could alter or modify the required chemical-physical properties of the drum material (chromium stainless steel, AISI 430), which is would be seriously affected by a thermal treatment like that deriving from a spot welding operation. On the other hand, this spot welding operation, which must involve a number as low as possible of regions of the drum, does not provide a connection having the desired strength and duration. In fact, the welding spots do not provide a structurally continuous connection and, moreover, as stated, could undesirably modify the drum material thereby favoring oxidation phenomena which would be very dangerous for the time stability of the performed connection.\nIn addition to the above mentioned drawbacks, to weld the front ring element or cover on the wedge elements of the drum is disadvantageous since it does not allow a proper centering of these elements. Actually, as the welding operation is performed, the ring or cover element can be displaced from a proper position centered with respect to the cylindric band on which they must be locked, thereby the end product may be affected by openings or undesired gaps between the connected portions.\nTo the foregoing it should be added that for welding the drum, additional welding materials and power are required, thereby increasing the end cost of the drim assembly.\nAccordingly, the aim of the present invention is to provide such a drum assembly for washing, drying and the like machines which has a mechanical strength greater than that of commercially available like devices, mainly at the level of the mutual connection regions coupling the mentioned wedge elements and front ring or cover element.\nWithin the scope of the above mentioned aim, a main object of the present invention is to provide a drum assembly of the above mentioned type which is effectively devoid of any unevenness like those which occur in the connection regions of the wedge elements and the front ring or cover element in prior drum assemblies.\nYet another object of the present invention is to provide a method for making the above mentioned drum assembly, which is specifically designed for reducing to a minimum the end cost of washing, drying and the like machine drum assemblies.\nAccording to one aspect of the present invention, the above mentioned aim and objects are achieved by an improved drum assembly for washing, drying and the like machines, of the type including a cylindric band provided with a plurality of wedge elements facing the inside of the drum assembly, a closing rear disc and a front ring element, characterized in that said drum assembly comprises connecting means for directly mechanically connecting said ring element and band, at said wedge elements of said band.\nAccording to further features of the inventive drum assembly, said connecting means comprise projecting or raised walls provided on that section of each said wedge element engaging said ring element, said projection portions being adapted to cooperate with corresponding projecting or raised portions of said ring element. Said projecting portions, in particular, comprise a wall formed on the end of the wedge elements facing the ring element and directed inwardly of said wedge elements, at a position substantially parallel to the surface of the ring element therewith said wall must be connected. In turn, said projecting portions of said ring element comprise a deformed or upset edge or rim portion provided at the connection section on each wedge element and having such a size and shape to be housed inside each said wedge element, above and beyond the mentioned projecting wall.\nAccording to further features of the drum assembly according to the present invention, said deformed edge or rim portion comprises a wall facing the inside of the drum assembly. After having performed the mechanical connection, the mentioned deformed edge or rim portion will be rearwardly turned and upset against said wall of the respective wedge element.\nAccording to a modified embodiment of the drum assembly of the invention, said direct connection mechanical means comprise a portion of said wedge elements axially projecting on the outside of said cylindric band, for cooperating with a corresponding window of the front ring element, to be locked inside the latter, at the contour thereof.\nThe method for making the above disclosed drum assembly, which also constitutes a main aspect of the present invention, is essentially characterized in that it comprises the step of cold anchoring said front ring or cover element and cylindric band, at the wedge elements of said band. The method is moreover characterized in that said anchoring or connecting step is performed by an upsetting operation which is preceded by a self-centered connection between said front ring or cover element and said cylindric band of the drum assembly.\nFurther features of the inventive method are clearly defined in the remaining dependent claims.\nWith respect to prior drum assemblies, the inventive drum assembly provides the advantage of allowing a mechanical direct connection of the wedge elements and front ring or cover element, i.e. a connection which does not comprise any welding spots favouring oxidation and sealing failure phenomena of the performed connection.\nThe mechanical sealing properties of the subject connection are moreover much more stable and efficient, than prior welding spot connections, since the subject connection is extended to the overall mutual coupling section of the wedge elements and corresponding portion of the front ring or cover element, and not to only a portion thereof, as in a spot welded connection.\nThe cold connection performed by upsetting a projecting portion of said front ring or cover element on the contour of a corresponding wedge element, furthermore, provides a self-centering effect. Thus, the drum assembling operation can be carried out a fully automatic manner due to the mutual engagement of the portions being connected, thereby providing a well finished final product."}
{"text": "1. Technical Field\nThe present invention relates to an electrostatic charge image developing toner and a manufacturing method thereof, an electrostatic charge image developer, a toner cartridge, a process cartridge, an image forming apparatus, and an image forming method.\n2. Related Art\nIn recent years, image forming apparatuses such as a printer and a copying machine have become widespread, and technologies related to various elements making up the image forming apparatuses have become widespread. Among the image forming apparatuses, in an image forming apparatus adopting an electrophotographic type, frequently, a photoreceptor including a photoreceptor (image holding member) is charged using a charging device, and an electrostatic charge image having a potential different from an ambient potential is formed on the charged photoreceptor to form a pattern that is desired to be printed. The electrostatic charge image that is formed in this manner is developed with toner and is ultimately transferred onto a transfer medium such as a recording paper."}
{"text": "There exists a need for adaptive optics systems which can correct for light distortions which occur in the medium of transmission. Conventional adaptive optics (“AO”) systems measure the characteristics of a lens, and correct for distortions using a computer-controlled deformable mirror, where the device which measures the distortions in the incoming wavefront of light is known as a wavefront sensor.\nFIG. 1 depicts conventional AO system 100, which is an example of a closed loop AO system commonly used in the astronomical community. Light beam 101 emanating from nominal point source 102, a star, is directed through reflective and refractive elements within telescope 104, toward fast steering mirror 105. Fast steering mirror 105 directs light beam 101 into deformable mirror 106 and the wavefront sensor 112, where wavefront errors are corrected, as described more fully below.\nLight beam 101 enters beam splitter 107, where it is split into two beams. Beam 109 enters science camera 110, where the light beam is viewed by the user of conventional AO system 100. Light beam 111 is directed into Shack-Hartmann wavefront sensor 112, which measures phase errors representing the derivative or slope of the wavefront in order to correct beam train or path length aberrations. These errors are communicated to wavefront computer 114, which calculates a reconstruction matrix to determine the phase or integral of the slope and passes these values as correction voltages to be applied to fast steering mirror 105 as X/Y tip/tilt values to stabilize the image, and phase values Φij (representing the phase values of an array of i,j pixels) to the deformable mirror 106, which restores the image sharpness lost to atmospheric turbulence. Wavefront computer 114 transmits these X/Y tip/tilt and Φij phase correction voltages to mirror drivers 115 and 116. Mirror drivers 115 and 116 apply the correction voltages to fast steering mirror 105 and deformable mirror 106, respectively, which apply the X/Y tip/tilt correction and Φij phase correction to the incoming wavefront from the target or star 102.\nFIG. 2 depicts the wavefront sensing technique used by Shack-Hartmann wavefront sensor 112. As shown, this approach is completely geometric in nature, and has no dependence on the coherence of the sensed optical beam. The wavefront of light beam 111 is broken into an array of spatial samples, called subapertures of the primary aperture, by two dimensional array of lenslets 202. The subaperture sampled by each lenslet is brought to a focus at a known distance F behind each array, on two-dimensional detector array 204. The lateral position of the focal spot depends on the local tilt of the incoming wavefront. The location of all the subaperture spot positions and corresponding spot deviations ΔX and ΔY are used to measure the gradient of the incoming wavefront. A two-dimensional integration process, known as ‘reconstruction,’ is then used to estimate the shape of the original wavefront, and thus derive the correction signals for the deformable mirror.\nTo its disadvantage, however, conventional Shack-Hartmann AO sensor 112 requires massive reconstructors, and the sensor is not highly scalable. Other conventional methods, such as heterodyne direct phase methods or active imaging, have high laser power requirements, or have a high minimum pixel count. Accordingly, it is desirable to provide for wavefront sensing which overcomes the deficiencies of conventional approaches. In particular, it is desirable to provide for the closed loop correction of phase aberrations, using a black fringe wavefront sensor, to avoid the use of a reconstructor."}
{"text": "Hindered amine light stabilizer (“HALS”) molecules may be added to a polymer in order to protect the polymer from radiation damage, such as ultraviolet (UV) degradation of the polymer. HALS molecules are believed to provide protection from radiation damage by terminating photo-oxidation processes in polymers through chemical reaction with free radical and peroxide intermediates. A common approach to render a polymer flame retardant is by incorporation of additives such as halogenated (e.g., brominated) materials. In some cases, brominated flame retardant additives may release bromine radicals that may react directly with the HALS molecules or may abstract a hydrogen from the polymer matrix and deactivate the HALS molecules through an acid-base reaction. The result is loss of light stabilization and rapid UV degradation of the unprotected polymer."}
{"text": "The term powertrain, or powerplant, describes the main components that generate power and deliver power to the road surface, water, or air. In a wider sense, the powertrain includes components used to transform a stored (chemical, solar, nuclear, kinetic, potential, etc.) energy into kinetic energy for propulsion purposes.\nThe powertrain of the current and upcoming electric vehicles (EVs) and plug-in hybrid electric vehicles (PHEVs), cumulatively called plug-in electric vehicles (PEVs), is composed of a high-power battery pack, as the Energy Storage Sub-System (ESS), and an Inverter followed by a Propulsion Machine for electric propulsion.\nElectric vehicles are vehicles propelled by electricity as opposed to the conventional vehicles which operate by liquid or gaseous fuels. Electric vehicles may use a combination of different energy sources such as, for example, fuel cells (FCs), batteries, and super-capacitors (or ultra-capacitors) to power an electric drive system as shown in FIG. 1. Electric vehicles (EVs) operate with high energy conversion efficiency, produce a lower level of exhaust emissions, and lower level of acoustic noise and vibration than conventional vehicles. The electricity needed for electric vehicles' operation may be produced either external the vehicle and stored in the battery, or produced onboard with the help of fuel cells.\nA plug-in electric vehicle (PEV) is any motor vehicle that can be recharged from an external source of electricity, such as wall sockets. The electricity stored in the rechargeable battery packs drives or contributes to drive the wheels. PEV is a superset of electric vehicles that includes all-electric or battery electric vehicles (BEVs), plug-in hybrid vehicles (PHEVs), and electric vehicle conversions of hybrid electric vehicles as well as conventional internal combustion engine vehicles.\nIn electric vehicles, the main energy source is assisted by one or more energy storage devices. A combination of batteries and super-capacitors are often used as Energy Storage Sources which can be connected to the fuel cell stack in a number of ways. The voltage characteristics of the two devices must match perfectly, and only a fraction of the range of operation of the Energy Storage devices can be utilized. For example, in a fuel cell/battery configuration, the fuel cell must provide substantially constant power all the time due to the fixed voltage of the battery, and in a battery/super-capacitor configuration, only a fraction of the energy exchange capability of the super-capacitor (also referred to herein as ultra-capacitor) can be used.\nDC-to-DC converters can be used to interface the elements in the electric powertrain by boosting or chopping the voltage levels as required by the load in different regimes of the EV operations. By introducing the DC-to-DC converters, one can select the voltage variations of the devices, and the power of each device can be controlled.\nDifferent configurations of EV power supplies dictates that at least one DC-to-DC converter is necessary to interface the fuel cell, the battery, or the super-capacitor's module and the DC Link. The DC-to-DC converter structure has to be reliable and lightweight, small in volume, operate with high efficiency, and low electromagnetic interference, and low current/voltage ripple.\nA DC-to-DC converter is a category of power converters having an electric circuit which converts a source of direct current (DC) from one voltage level to another, by storing the input energy temporary and subsequently releasing the energy to the output at a different voltage level.\nThe storage of the input energy may be in either magnetic field storage components (inductors, transformers) or electric field storage components (capacitors). DC-to-DC converters can be designed to transfer power in only one direction, from the input to the output. However many all DC-to-DC converter topologies can be made bi-directional which can move power in either direction which is useful in applications requiring regenerative braking.\nThe amount of power flow between the input and the output of the DC-to-DC converter can be controlled by adjusting the duty cycle (which is identified as a ratio of on/off time of a switch in the DC-to-DC converter). Usually, this is done to control the output voltage, the input current, the output current, or to maintain a constant power.\nIt is desirable to use DC-to-DC converters with minimized switching losses, lower inductor current ripple and low switch voltage ratings, as well as reduced output voltage ripples, and enhanced reliability and autonomity of the operation in plug-in electric vehicles.\nIt is also desirable to provide various configurations of power electronic interfaces between the Energy Storage Sub-Systems and the Propulsion Machines (including electric traction system), in the plug-in electric vehicles, with efficient dynamic splitting of power output from the Energy Storage Sub-Systems to power the Propulsion Machines in accordance with their operational power requirements in a dynamic fashion in order to attain an extended ESS usability and lifetime and to provide the power flow control flexibility."}
{"text": "Prior art tiltable cargo racks for supporting large pre-cast concrete structures had a longitudinal arm rotatable from a horizontal loading position to an inclined transport position. An example of such a rack is described in U.S. Pat. No. 5,947,665 to Baur et al., FIGS. 1-4 and the corresponding text of U.S. Pat. No. 5,947,665 are incorporated herein by reference in their entirety. In such prior cargo racks, a plurality of upwardly extending members of specified length was required for simultaneously contacting and supporting respective corners of the structure. Adequate support was not provided in the event that the pre-cast concrete structure had one or more corners at a height differing from the length of one or more of the corresponding upwardly extending members of the rack. Also, adequate support was not provided to prevent unintentional tilting of the rack during cargo loading.\nAccordingly, it is desirable to provide a tiltable cargo rack that provides adequate support for the cargo without requiring a plurality of upwardly extending members of specified length that simultaneously contact and support respective corners of the cargo."}
{"text": "Resistance change memories such as a magnetic random access memory (hereinafter referred to as an MRAM) that uses magnetoresistive effects of a ferromagnetic substance have been recently drawing more attention as next-generation solid-state nonvolatile memories capable of high-speed reading and writing, having high capacities, and capable of low-power-consumption operation. In particular, a magnetoresistive element having a magnetic tunnel junction used in the MRAM is also referred to as a magnetic tunnel junction (MTJ) element, and has been dawning attention since it had been proved to show a high magnetoresistance rate. The MTJ element has a three-layer stacked layer structure including a storage layer variable in magnetization direction, an insulator layer, and a reference layer which faces the storage layer and which maintains a predetermined magnetization direction."}
{"text": "This section is intended to provide a background or context to the invention recited in the claims. The description herein may include concepts that could be pursued, but are not necessarily ones that have been previously conceived or pursued. Therefore, unless otherwise indicated herein, what is described in this section is not prior art to the description and claims in this application and is not admitted to be prior art by inclusion in this section.\nTemperature-controlled storage devices (e.g., a refrigerator, freezer, refrigerated merchandiser, display case, etc.) are used in a wide variety of commercial, institutional, and residential applications for storing and/or displaying refrigerated or frozen objects. Many temperature-controlled storage devices have a display case door (e.g., a door with an insulated glass panel) through which objects within the temperature-controlled storage device can be viewed. Traditional display case doors allow a customer in a supermarket or other similar facility to see the objects within the temperature-controlled storage device while shopping, but do not otherwise facilitate product presentation or enhance the shopping experience."}
{"text": "1 Field of the Invention\nThis invention relates to an improved system for integrating automatic outbound dialer functions with automatic call distribution functions, and more particularly to a system that balances the number of agents assigned to these functions while maintaining the inbound call waiting time within prescribed limits.\n2 Description of the Prior Art\nMany call centers provide both an automatic inbound telephone call distribution function and an automatic outbound call function. As will be appreciated by those skilled in the art, there have been proposals in the prior art to link inbound automatic call distribution and automatic outbound calling in order to improve staffing efficiency. Inbound call distribution typically has peaks and valleys in its load, since the demand is generated by outside callers. By linking inbound and outbound call functions, outbound agents can be switched to inbound duty during peak inbound demand periods and switched back to outbound duty during slack periods in inbound demand, thus improving overall staffing efficiency.\nIn a typical prior art system, an inbound performance parameter is monitored based upon statistics tracked by the inbound call distributor; for example, number of calls in the inbound queue or average time to answer an inbound call. Target values, and upper and lower thresholds are established for the inbound performance parameter; for example a five second target to respond to incoming calls with an upper threshold of seven seconds and a lower threshold of three seconds. The assumptions of the thresholds is; if upper threshold is exceeded, the performance is unacceptable and more inbound agents are needed; and (b) if the lower threshold is exceeded, the inbound function is considered to be overstaffed and overall efficiency would be improved by transferring inbound agents to outbound operations.\nIn the prior art, so called blend agents are typically transferred, as they become available, when the upper or lower threshold values are exceeded. Agents continue to be transferred until the performance is between the upper and lower threshold limits, when the transfer stops. U.S. Pat. No. 5,425,093, assigned to the assignee of this application and incorporated herein by reference, discloses the concept of providing hysteresis to provide added stabilization to agent transfer. With hysteresis, a higher threshold value is used to trigger the transfer of agents as the upper range of acceptable performance is being exceeded than the threshold value used to stop transfer as the performance moves toward the acceptable range. Similarly for the lower limit; the trigger value to start transfer is lower as measured performance leaves the range than to stop transfer as the performance moves toward the range. However, this and other prior art call center balance algorithms for determining the transfer of agents between inbound and outbound functions operate in response to current inbound performance, and have not proven to be altogether satisfactory in operation."}
{"text": "This invention relates generally to the field of telephony, and more particularly to an improved form of surge protection unit particularly adapted to be installed upon existing subscriber lines in such locations as motels, small apartment or office buildings and the like without a complicated installation procedure.\nIn such locations, the existing subscriber lines which are usually twenty-five or less in number are often connected to the telephone company lines through a small connector block of elongated configuration commonly known as a Western Electric Type 66 block, which block has been used in the industry for many years. Blocks of this type are characterized in the provision of an exposed surface mounting for parallel rows of quick clip connectors which are connected through wire wrap means to the telephone company lines, and through the expandable clip portions of the quick clip connectors to the subscriber lines.\nWith the recent emergence of subscriber owned equipment, and a lessening of maintenance responsibility on the part of the telephone company, it is desirable to provide a surge protection unit at or near the subscriber premises. In the case of buildings having up to twenty-five individual circuits, the use of the Type 66 block is commonly found in multiple installations. On occasion, multiples of this type of block are used in physically juxtaposed locations."}
{"text": "Conventionally, there is a multi function peripheral (hereinafter referred to as an “MFP”) and an image forming apparatus such as a printer and the like. The image forming apparatus is provided with a fixing device. The fixing device heats a fixing belt through an electromagnetic induction heating method (hereinafter referred to as an “IH method”) to fix a toner image on an image receiving medium through the heat of the fixing belt. The fixing belt has a heating layer which generates heat through induction current. The fixing device reduces the heat capacity of the fixing belt to reduce the warming up time and the like. The fixing device includes a magnetic shunt alloy layer for compensating the deficiency of the calorific value of the fixing belt. The magnetic characteristic of the magnetic shunt alloy layer varies according to the temperature. The magnetic shunt alloy layer is transformed from a ferromagnetic body into a paramagnetic body at a curie point. The magnetic shunt alloy layer generates heat itself. The magnetic shunt alloy layer loses its magnetism and does not generate heat any more at the curie point. The calorific value of the fixing belt may be excessively increased due to the heat generated by the magnetic shunt alloy layer itself, which may lead to a problem that the heating efficiency of the fixing belt cannot be improved."}
{"text": "The field of the invention pertains to electroacoustic loudspeakers for home and concert stereophonic music reproduction. In particular, the invention is directed to modifications of such loudspeakers to reduce the effect of the loudspeaker's electroacoustic configuration on the sound signal and thereby to produce a more accurate reproduction of music.\nRecently, a variety of loudspeakers have been introduced that are directed to creation of ambient sound fields in addition to the direct sound fields generated by drivers directed at the expected listener position. Such loudspeakers incorporate rear facing drivers and upward or downward facing drivers.\nLoudspeakers such as the \"TIME WINDOW\" from DCM CORPORATION of Ann Arbor, Mich. employ a prismatic or dihedral front with drivers in both panels. The drivers in the panels facing the listener provide the direct sound field and the drivers in the panels facing away from the listener provide the ambient or reflected sound field.\nReference is made to applicant's U.S. Pat. No. 4,315,102 which discloses improvements in electric crossover circuitry and methods to adjust loudspeakers for flatter amplitude and phase response within and beyond the pass bands of the electroacoustic networks. Reference is also made to applicant's pending U.S. application Ser. No. 146,254 which discloses improvements in crossovers incorporating delay lines, coaxial drivers and methods to adjust loudspeakers for flatter amplitude and phase response within and beyond the pass bands of the electroacoustic networks. The above applicant's disclosures are incorporated herein by reference.\nU.S. Pat. No. 2,053,364 and U.S. Pat. No. 2,259,907 illustrate early forms of coaxial high and low frequency driver combinations for incorporation in a loudspeaker. The high frequency driver is located in front of the low frequency driver which unfortunately accentuates the difference in time of arrival at the listener's ear of sounds generated simultaneously by the two drivers of the coaxial driver.\nThe delay lines disclosed in applicant's above noted pending application are directed toward overcoming some of the distortions in the acoustic signals reaching the ear that arise from the use of coaxial drivers. The disclosure below comprises further improvements to the direct and ambient time/space coherent sound approach pioneered by applicant with the original \"TIME WINDOW\" loudspeakers and with other loudspeaker products of DCM CORPORATION."}
{"text": "In heat assisted magnetic recording, information bits are recorded on a data storage medium at elevated temperatures, and the data bit dimensions can be determined by the dimensions of the heated area in the storage medium or the dimensions of an area of the storage medium that is subjected to a magnetic field. In one approach, a beam of light is condensed to a small optical spot onto the storage medium to heat a portion of the medium and reduce the magnetic coercivity of the heated portion. Data is then written to the reduced coercivity region.\nOne example of a recording head for use in heat assisted magnetic recording generally includes a write pole and a return pole magnetically coupled to each other through a yoke or pedestal, and a waveguide for focusing light onto the storage medium. The waveguide is positioned between the write pole and the return pole. Light from a laser is coupled into the waveguide and transmitted through the waveguide to produce a tightly focused laser spot on the storage media. The power in this spot determines the width of the recorded track and the quality of the recorded transitions.\nIt would be desirable to have a method to measure or monitor the optical power once the head is assembled in a drive."}
{"text": "A network generally comprises equipments to be managed and a device for managing those equipments.\nNetwork management is based on the management device and the equipments to be managed exchanging information. Thus the management device can send commands to one or more network equipments, in particular to request modification of configuration parameters, or merely to consult certain configuration parameters of the network equipments. The network equipments prompted in this way generally modify the parameters on receiving a modification request or send requested information on receiving an information request. The network equipments are also adapted to send information spontaneously to the management device, such as in particular information relating to alarms.\nThe data exchanged in this way must be comprehensible to the sender and to the receiver. To this end, a network conventionally uses a management information base (MIB) reference model of such management information.\nAn MIB specifies a management information model in the form of objects to be managed in relation to parameters to be consulted, parameters to be modified, or alarms to be sent. These objects are specified by the International Standardization Organization (ISO) in the form of an object record tree. Thus a reference is assigned to each object specified in the MIB.\nAs a function of the type of management required, it is possible to refer to only a portion of the MIB tree, below called the “MIB portion”.\nThe MIB contains individual scalar objects that can be instantiated once only, and groups of objects that can be instantiated a varying number of times, these groups being conventionally represented in the form of tables in which each row represents one instance of the group of objects. Each row is then identified by an index.\nIn a management network, the equipments and the management device generally know the MIB portion in accordance with which they manage their own management data. Thus a network is conventionally managed on the basis of a model of the above type, with knowledge thereof being shared by the management device and the network equipments to be managed.\nA standardized protocol such as the simple network management protocol (SNMP) is generally used to send this information between the equipments and the management device. This kind of protocol is suited to exchanging management data relating to objects specified in the MIB, but it occupies considerable transmission bandwidth in the network. About 75% of the size of each message corresponds to ‘Abstract Syntax Notation 1’/‘Basic encoding rule’ (ASN.1/BER) references."}
{"text": "Bulk Acoustic Wave (BAW) filter and Bulk Acoustic Resonator (BAR) are gaining more popularly for their performance benefits and are being utilized in the design of today's cutting-edge mobile devices and systems. However, due to manufacturing complexities, Bulk Acoustic Wave (BAW) filter and Bulk Acoustic Resonator (BAR) are fabricated as standalone devices. That is, the Bulk Acoustic Wave (BAW) filter and Bulk Acoustic Resonator (BAR) are not provided as integrated structures with other CMOS, BiCMOS, SiGe HBT, and/or passive devices, thus leading to higher manufacturing costs, and increased fabrication processing.\nAccordingly, there exists a need in the art to overcome the deficiencies and limitations described hereinabove."}
{"text": "Disposable sheet style dispensers are well-known in the art for dispensing individual folded sheet products such as facial tissues, hand sheets, wet wipes, or the like. In general, disposable sheet dispensers typically include a container and a stack or clip of pre-folded, interfolded sheets disposed within the container. The sheets may be C-folded, V-folded, or flat so that once the top sheet in the clip is withdrawn, the underlying sheet is individually presented for subsequent use.\nOne of the common problems among disposable sheet dispensers involves the issue of xe2x80x9cdouble pullxe2x80x9d. This occurs when more than one sheet comes out of the container when the leading sheet is withdrawn. Additionally, the problem of xe2x80x9cstreamingxe2x80x9d can occur in disposable sheet dispensers. Streaming occurs when the user pulls the first sheet out, and subsequent sheets are also withdrawn, with the separation of the following ones never occurring.\nTo overcome some of these problems, disposable sheet dispensers have been designed where the dispensing slot is a slit. The use of slits has worked well in eliminating some of the above problems. Narrow slits are also well-suited to presenting tissue sheets for subsequent removal. In some applications, however, it is sometimes difficult to remove the initial sheet from the stack when a slit is present. For instance, the slit can interfere with the ability of a user to grab the leading edge.\nIn order to overcome some of these problems, in U.S. Pat. No. 4,574,952 to Masui, a box containing facial tissues is disclosed in which a tape or strings are attached to the undersurface of the top sheet of the box and, in turn, attached to the upper most of the facial tissues. In this manner, when the box is opened, and the top sheet is removed along a perforated line, the upper most facial tissue is automatically removed from the container along with the top sheet.\nAlthough the top facial tissue is automatically removed from the box when the box is opened, the tissue may not be usable since it is attached to the removable top sheet of the box. Further, many times tissue boxes are opened for display purposes where a tissue is not initially needed. In these instances, removal of the first tissue by automatically opening the box can lead to waste of the leading sheet. In view of these drawbacks, the present invention is directed to further improvements in interfolded sheet dispensers.\nIn general, the present invention is directed to a dispenser for dispensing interfolded disposable sheets. A stack of interfolded disposable sheets are housed within a container. The sheets are arranged to be withdrawn from the container one after another. The sheets can be various paper products, including facial tissues, paper towels, industrial wipers, laboratory wipers, wet wipes, and the like. The dispensing container includes a removable panel that, once removed, uncovers an opening for withdrawing the sheets.\nIn accordance with the present invention, the removable panel is attached to a first sheet in the stack of interfolded sheets. In particular, the removable panel is attached to the first sheet such that when the removable panel is removed from the dispensing container, the first sheet is partially pulled through the opening of the dispensing container without being completely removed from the container.\nAn attachment device is used to attach the first sheet in the stack to the removable panel. The attachment device can be, for instance, a mechanical structure or an adhesive. In one embodiment, the attachment device is a strip of double-sided adhesive tape. For example, in one embodiment, the double-sided adhesive tape has a first side that is adhered to the removable panel and a second side that is adhered to the first sheet. In this embodiment, the first side can have a greater adhesive strength than the second side.\nIn one embodiment, the first sheet has a first end that is interfolded with another sheet in the stack. The second end of the sheet, however, is attached to the removable panel and is folded upon itself in an amount sufficient for the first sheet to be partially removed from the dispensing container when the removable panel is removed. In this manner, resistance to being removed from the container is increased as the first sheet is pulled through the opening by the removable panel.\nThe dispensing container can be made from various materials and is generally not critical to the present invention. For instance, the dispensing container can be made from paperboard or from a flexible polymer film. The dispensing container can be in the shape of a rectangular box, a square box, or in the form of any other suitable shape.\nOther features and aspects of the present invention are discussed in greater detail below."}
{"text": "An MR sensor detects magnetic field signals through the resistance changes of a read element, fabricated of a magnetic material, as a function of the strength and direction of magnetic flux being sensed by the read element. The conventional MR sensor, such as that used in the IBM \"Corsair\" disk drive, operates on the basis of the anisotropic magnetoresistive (AMR) effect in which a component of the read element resistance varies as the square of the cosine of the angle between the magnetization in the read element and the direction of sense current flow through the read element. Recorded data can be read from a magnetic medium because the external magnetic field from the recorded magnetic medium (the signal field) causes a change in the direction of magnetization in the read element, which in turn causes a change in resistance in the read element and a corresponding change in the sensed current or voltage.\nA different and more pronounced magnetoresistance, called giant magnetoresistance (GMR), has been observed in a variety of magnetic multilayered structures, the essential feature being at least two ferromagnetic metal layers separated by a nonferromagnetic metal layer. This GMR effect has been found in a variety of systems, such as Fe/Cr, Co/Cu, or Co/Ru multilayers exhibiting strong antiferromagnetic coupling of the ferromagnetic layers, as well as in essentially uncoupled layered structures in which the magnetization orientation in one of the two ferromagnetic layers is fixed or pinned. The physical origin is the same in all types of structures: the application of an external magnetic field causes a variation in the relative orientation of neighboring ferromagnetic layers. This in turn causes a change in the spin-dependent scattering of conduction electrons and thus the electrical resistance of the structure. The resistance of the structure thus changes as the relative alignment of the magnetizations of the ferromagnetic layers changes.\nA particularly useful application of GMR is a sandwich structure comprising two uncoupled ferromagnetic layers separated by a nonmagnetic metallic layer in which the magnetization of one of the ferromagnetic layers is pinned. The pinning may be achieved by depositing the layer onto an iron-manganese (Fe--Mn) layer to exchange couple the two layers. This results in a spin valve magnetoresistive (SVMR) sensor in which only the unpinned or free ferromagnetic layer is free to rotate in the presence of an external magnetic field. U.S. Pat. No. 5,206,590, filed Dec. 11, 1990, and assigned to IBM, discloses a basic spin valve MR sensor. U.S. Pat. No. 5,159,513, filed Feb. 8, 1991, and also assigned to IBM, discloses a spin valve MR sensor in which at least one of the ferromagnetic layers is of cobalt or a cobalt alloy, and in which the magnetizations of the two ferromagnetic layers are maintained substantially perpendicular to each other at zero externally applied magnetic field by exchange coupling of the pinned ferromagnetic layer to an antiferromagnetic layer.\nThe spin valve MR sensor that has the most linear response and the widest dynamic range is one in which the magnetization of the pinned ferromagnetic layer is parallel to the signal field and the magnetization of the free ferromagnetic layer is perpendicular to the signal field. In the case where the spin valve MR sensor is to be used in a horizontal magnetic recording disk drive, this means that the plane of the sensor is perpendicular to the disk surface with the magnetization of the pinned layer oriented perpendicular to and the magnetization of the free layer oriented parallel to the disk surface. One difficulty in achieving this magnetization orientation is caused by the dipole field generated by the pinned layer. The pinned layer has a net magnetic moment and thus essentially acts as a macroscopic dipole magnet whose field acts on the free layer. In spin valve MR sensors, where the read element is relatively short, the result of this magnetostatic coupling is that the magnetization direction in the free layer is not uniform. This causes portions of the sensor to saturate prematurely in the presence of the signal field, which limits the sensor's dynamic range and thus the recording density and overall performance of the magnetic recording system.\nWhat is needed is a spin valve MR sensor with an improved pinned ferromagnetic layer that has minimal magnetostatic coupling with the free ferromagnetic layer."}
{"text": "The present invention relates to the art of providing protection for active food ingredients, and, in particular, to protecting food ingredients by encapsulation.\nIt has been known in the art of food preparation to provide protection for different ingredients included in food compositions by means of protective coatings. Such protective systems have been employed for various reasons such as for protection of the active ingredient, both while on the shelf and during use, and for prolonged release in the oral cavity and/or systemically. As used herein, \"active ingredient\" can be an ingredient such as a sweetener, soluble dietary fiber, a flavoring agent, a bio-effecting ingredient, such as a medicinal drug or pharmaceutical agent, and a breath-freshening ingredient.\nProtective systems for active ingredients encounter unique problems when used in chewing gum compositions since gum compositions include a soluble portion which is consumed during chewing and an insoluble masticatory portion which remains in the oral cavity as a cud or bolus throughout chewing. The soluble portion is added by blending soluble bulking agents, flavorants, natural and synthetic sweeteners, etc. to the gum base usually in the presence of heat in order to effectuate mixing. Thus, a protective system for an active must be able to withstand the rigors of heat and shear forces. Furthermore, since the gum base includes basically insoluble ingredients such as resins, elastomers, etc., any protection for an active ingredient included in the base would have to be compatible to a certain degree with such components. Depending on the effect desired, the protected active ingredient can be incorporated into the base and/or added with the soluble components or added after compositing.\nDipeptide sweeteners such as L-aspartyl-L-phenyl-alanine methylester (aspartame), which have been widely recognized as an extremely good-tasting non-caloric sweetener for use in a variety of food products, have been found to be particularly troublesome, especially, when used in chewing gum products. Unfortunately, aspartame is extremely unstable in the presence of moisture and undergoes hydrolytic degradation and subsequent loss of sweetness. Elevated temperatures and specific pH ranges may accelerate the hydrolytic degradation. Additionally, aspartame is known to react with a number of flavorings and chemicals such as aldehydes and ketones. For example, aspartame loses its sweetness due to chemical reactions in the presence of flavor aldehyde-containing oils such as cinnamon. These flavor oils are used in food products and are popular in chewing gum and other confectionery products. These compositions also generally contain moisture and may be exposed to elevated temperatures during their shelf life. The result is the loss of aspartame and the corresponding loss of sweetness contributed thereby in the gum.\nFor example, U.S. Pat. No. 3,962,463 to Witzel discloses chewing gum containing on its surface tiny capsules of flavoring and/or flavors fixed on an edible substrate which will release the flavor upon mastication. The flavoring ingredient can be micro-encapsulated in gelatin, waxes, polyethylene and the like, and printed on the surface of the gum as an aqueous slurry. The encapsulation can be effected in a conventional manner by blending liquid flavoring with a concentrated aqueous solution of gelatin at a temperature below 25.degree. C. whereby a fine stable emulsion is formed. After treatment to impart moisture resistance, the emulsion can be spray dried while still cool to produce a fine free-flowing powder each particle of which consists of a core of flavoring surrounded by a dry gelatin wall. The gelatin micro-encapsulated flavoring ingredient can be printed on the gum by either direct or off-set Gravure printing. As can be seen, the provision of a protection system according to the above-disclosure is a rather complex procedure which requires an additional manufacturing step in order to be added to the gum product.\nOther disclosures show the use of protective encapsulation of artificial sweeteners such as aspartame, saccharin, etc. For example, U.S. Pat. Nos. 4,122,195 and 4,139,639 disclose encapsulation of aspartame in Capsul dextrin and gum arabic. However, such encapsulants have been found to be of only limited effectiveness in preventing degradation of moisture sensitive aspartame since these encapsulants are hydrophilic and moisture-permeable.\nU.S. Pat. No. 4,384,004 to Cea, et al. discloses preparation of an artificial sweetener, i.e., APM, in an encapsulated form with at least one layer of a coating material selected from the group consisting of cellulose, cellulose derivatives, starches, carbohydrates, gums, polyolefins, polyesters, waxes, vinyl polymers, gelatin, zein, and mixtures thereof in a ratio of coating material to APM which does not exceed 1:1. This coating material is applied to the APM while it is in solid particle form at temperatures below the decomposition temperature of APM, e.g., up to about 200.degree., and preferably about 100.degree. to 115.degree. F. The application of the Cea, et al. protective system requires a rather elaborate energy-intensive coating procedure, e.g., fluidized bed coating process, in which a strong upward stream or air current must be provided. This results from, among other things, the requirement of providing the protective coating material in a solvent.\nEPA No. 81110320.0, published June 16, 1982 (Publication No. 0053844), to Ajinomoto Co., Inc., discloses a stabilized dipeptide-based sweetening composition comprising (a) from 20 to 60% by weight of solid fat, (b) from 10 to 30% by weight emulsifier, (c) from 10 to 30% by weight polysaccharide and (d) not more than 30% by weight of dipeptide sweetener. The compositions are prepared by heating the mixture of the ingredients, cooling, and pulverizing to obtain powder of granules of the composition to obtain a ASTM mesh size of smaller than 12. Spray drying the mixture is also disclosed.\nU.S. Pat. No. 4,105,801 to Degliotti, discloses a confectionary comprising a core portion and a shell adheringly enveloping the core portion, whereby the shell is formed by an intimate mixture of microcrystals of xylitol with a solid fatty substance in a proportion of 0.5 to 15 parts by weight of fatty substance to each 100 parts by weight of xylitol. The fatty substance is preferably a mono-, di- or triglyceride having a melting range of between 20.degree. and 60.degree. C.\nWhile it would seem that hydrophobic encapsulating agents provide better impermeability and gradual release characteristics than hydrophilic coatings, it is not known whether any hydrophobic coatings have been successfully employed in the past for chewing gum ingredients. Most hydrophobic materials which could be used as encapsulants, such as a low molecular weight polyvinyl acetate, waxes and fats, are dissolved in the chewing base when they are mixed into the heated gum mass during the gum manufacturing process.\nOther hydrophobic materials such as high molecular weight polyvinyl acetate and styrene butadiene rubber are substantially insoluble in food grade solvents which are required in encapsulating processes. However, U.S. Pat. No. 3,826,847 to Ogawa, et al. discloses encapsulation of seasonings including sugar by use of high molecular weight polymers such as high molecular weight polyvinyl acetate; e.g., having a degree of polymerization of 400. Ogawa, et al. describe combining the seasoning by dissolving the polyvinyl ester in a solvent such as ethanol, ethyl acetate and the like to obtain a 2-30% by weight solution to which the seasoning is subsequently added with agitation in an amount at 1-20 times the polyvinyl ester content to form a homogenous dispersion. Then a liquid which is miscible with the solution but immiscible with the polyvinyl ester, such as ether, hexane, or the like, is slowly added to the dispersion so that polyvinyl ester granules integrated with the seasoning are separated out in the form of non-adhesive particles for incoporation into chewing gum base. The particle size can be less than 20 mesh, and preferably are less than 48 mesh. In order to produce the particulate, however, several decantations are required as well as a drying procedure, such as energy-intensive suction.\nYet another protective system for encapsulation has been disclosed in PCT/US No. 84,00108 including a shellac encapsulant, which is hydrophobic and insoluble in the gum base yet soluble in food-grade solvents, to provide a substantially impermeable coating for active ingredients, such as aspartame. The method of encapsulation can be effected by coating with a shellac-containing solution primarily by a fluidized bed coating method wherein particles of the active ingredient are suspended in a stream of pressurized air and sprayed with a solution of the encapsulating agent. The particles can also be coated in a roller bed coating method wherein a suspension of the active agent and its solvent is deposited on a heated, rotating drum, the solvent being evaporated by the heat leaving coated ingredient particles which are scraped from the roller. In each of these methods, and, generally, all methods of encapsulating which includes a solvent, an energy-intensive step to drive off the solvent is required before use in the end food product.\nIt is therefore, an object of the present invention to provide an active ingredient protection system which can be easily included in a food composition such as a chewing gum composition without an additional step in the manufacturing process.\nAnother object of the present invention is to provide an encapsulation system which is highly effective in preventing deterioration from moisture.\nA further object of the invention is to provide an encapsulation system which can be applied without the use of a solvent system.\nStill another object of the present invention is to provide an encapsulated active ingredient with a highly-controlled release characteristic for use in a food product.\nYet another object of the present invention is to provide a protective system for artificial sweeteners which deteriorate in the present of moisture for use in a chewing gum composition.\nAnother object of the invention is to provide a chewing gum composition with a highly-controlled active ingredient protection and release system.\nOther and further objects will become apparent to those skilled in the art in view of the disclosure set forth herein."}
{"text": "1. Field of the Invention\nThe present invention relates to an organopolysiloxane gel composition that yields a cured product via an addition reaction of an organopolysiloxane, and relates particularly to an organopolysiloxane gel composition that yields a cured product that can be used favorably for protecting ICs or hybrid ICs and for encapsulating power modules and the like.\n2. Description of the Prior Art\nSilicone gels and silicone potting materials are addition reaction-curable organopolysiloxane compositions which comprise an organohydrogenpolysiloxane having hydrogen atoms bonded to silicon atoms (namely, SiH groups), an organopolysiloxane having alkenyl groups such as vinyl groups bonded to silicon atoms, and a platinum-based catalyst, and which yield a cured product via an addition reaction of the SiH groups to the alkenyl groups such as vinyl groups. These silicone gels and silicone potting materials exhibit excellent levels of heat resistance, weather resistance, oil resistance, cold resistance, and electrical insulation, and also exhibit low elastic modulus values and low levels of stress, and as a result, they are widely used in the protection of vehicle-mounted electronic components and household electronic appliance components. In recent years, with the ongoing trend towards smaller and lighter electronic componentry, the density of IC components, and the density of the wiring such as wire bonding used to connect these IC components to substrates have continued to increase. Because conventional silicone gels are of comparatively low viscosity they exhibit favorable flow properties, but the cured compositions lack satisfactory strength and elongation properties, and are weak relative to large positional displacements, with the gel cured products sometimes rupturing. In contrast, it is known that compositions that yield cured products with satisfactory strength relative to displacement suffer from higher viscosity, with a corresponding deterioration in the flow properties during the encapsulation step.\nOrganohydrogenpolysiloxanes having SiH groups that also include branching points are disclosed in Japanese Patent Publication No. 2,510,577 and Japanese Patent Publication No. 2,849,027. However, the cured products of the gel compositions disclosed in these publications deform readily and are prone to rupture when subjected to a large displacement."}
{"text": "Scoliosis is a general term for the sideways (lateral) curving of the spine, usually in the thoracic or thoracolumbar region. Scoliosis is commonly broken up into different treatment groups, Adolescent Idiopathic Scoliosis, Early Onset Scoliosis and Adult Scoliosis.\nAdolescent Idiopathic Scoliosis (AIS) typically affects children between ages 10 and 16, and becomes most severe during growth spurts that occur as the body is developing. One to two percent of children between ages 10 and 16 have some amount of scoliosis. Of every 1000 children, two to five develop curves that are serious enough to require treatment. The degree of scoliosis is typically described by the Cobb angle, which is determined, usually from x-ray images, by taking the most tilted vertebrae above and below the apex of the curved portion and measuring the angle between intersecting lines drawn perpendicular to the top of the top vertebrae and the bottom of the bottom. The term idiopathic refers to the fact that the exact cause of this curvature is unknown. Some have speculated that scoliosis occurs when, during rapid growth phases, the ligamentum flavum of the spine is too tight and hinders symmetric growth of the spine. For example, as the anterior portion of the spine elongates faster than the posterior portion, the thoracic spine begins to straighten, until it curves laterally, often with an accompanying rotation. In more severe cases, this rotation actually creates a noticeable deformity, wherein one shoulder is lower than the other. Currently, many school districts perform external visual assessment of spines, for example in all fifth grade students. For those students in whom an “S” shape or “C” shape is identified, instead of an “I” shape, a recommendation is given to have the spine examined by a physician, and commonly followed-up with periodic spinal x-rays.\nTypically, patients with a Cobb angle of 20° or less are not treated, but are continually followed up, often with subsequent x-rays. Patients with a Cobb angle of 40° or greater are usually recommended for fusion surgery. It should be noted that many patients do not receive this spinal assessment, for numerous reasons. Many school districts do not perform this assessment, and many children do not regularly visit a physician, so often, the curve progresses rapidly and severely. There is a large population of grown adults with untreated scoliosis, in extreme cases with a Cobb angle as high as or greater than 90°. Many of these adults, though, do not have pain associated with this deformity, and live relatively normal lives, though oftentimes with restricted mobility and motion. In AIS, the ratio of females to males for curves under 10° is about one to one, however, at angles above 30°, females outnumber males by as much as eight to one. Fusion surgery can be performed on the AIS patients or on adult scoliosis patients. In a typical posterior fusion surgery, an incision is made down the length of the back and Titanium or stainless steel straightening rods are placed along the curved portion. These rods are typically secured to the vertebral bodies, for example with hooks or bone screws, or more specifically pedicle screws, in a manner that allows the spine to be straightened. Usually, at the section desired for fusion, the intervertebral disks are removed and bone graft material is placed to create the fusion. If this is autologous material, the bone is harvested from a hip via a separate incision.\nAlternatively, the fusion surgery may be performed anteriorly. A lateral and anterior incision is made for access. Usually, one of the lungs is deflated in order to allow access to the spine from this anterior approach. In a less-invasive version of the anterior procedure, instead of the single long incision, approximately five incisions, each about three to four cm long are made in several of the intercostal spaces (between the ribs) on one side of the patient. In one version of this minimally invasive surgery, tethers and bone screws are placed and are secured to the vertebra on the anterior convex portion of the curve. Currently, clinical trials are being performed which use staples in place of the tether/screw combination. One advantage of this surgery in comparison with the posterior approach is that the scars from the incisions are not as dramatic, though they are still located in a visible area, when a bathing suit, for example, is worn. The staples have had some difficulty in the clinical trials. The staples tend to pull out of the bone when a critical stress level is reached.\nIn some cases, after surgery, the patient will wear a protective brace for a few months as the fusing process occurs. Once the patient reaches spinal maturity, it is difficult to remove the rods and associated hardware in a subsequent surgery, because the fusion of the vertebra usually incorporates the rods themselves. Standard practice is to leave this implant in for life. With either of these two surgical methods, after fusion, the patient's spine is now straight, but depending on how many vertebra were fused, there are often limitations in the degree of flexibility, both in bending and twisting. As these fused patients mature, the fused section can impart large stresses on the adjacent non-fused vertebra, and often, other problems including pain can occur in these areas, sometimes necessitating further surgery. This tends to be in the lumbar portion of the spine that is prone to problems in aging patients. Many physicians are now interested in fusionless surgery for scoliosis, which may be able to eliminate some of the drawbacks of fusion.\nOne group of patients in which the spine is especially dynamic is the subset known as Early Onset Scoliosis (EOS), which typically occurs in children before the age of five, and more often in boys than in girls. This is a more rare condition, occurring in only about one or two out of 10,000 children, but can be severe, sometimes affecting the normal development of organs. Because of the fact that the spines of these children will still grow a large amount after treatment, non-fusion distraction devices known as growing rods and a device known as the VEPTR—Vertical Expandable Prosthetic Titanium Rib (“Titanium Rib”) have been developed. These devices are typically adjusted approximately every six months, to match the child's growth, until the child is at least eight years old, sometimes until they are 15 years old. Each adjustment requires a surgical incision to access the adjustable portion of the device. Because the patients may receive the device at an age as early as six months old, this treatment requires a large number of surgeries. Because of the multiple surgeries, these patients have a rather high preponderance of infection.\nReturning to the AIS patients, the treatment methodology for those with a Cobb angle between 20° and 40° is quite controversial. Many physicians proscribe a brace (for example, the Boston Brace), that the patient must wear on their body and under their clothes 18 to 23 hours a day until they become skeletally mature, for example to age 16. Because these patients are all passing through their socially demanding adolescent years, it is quite a serious prospect to be forced with the choice of either wearing a somewhat bulky brace that covers most of the upper body, having fusion surgery that may leave large scars and also limit motion, or doing nothing and running the risk of becoming disfigured and possibly disabled. It is commonly known that many patients have at times hidden their braces, for example, in a bush outside of school, in order to escape any related embarrassment. The patient compliance with brace wearing has been so problematic that there have been special braces constructed which sense the body of the patient, and keep track of the amount of time per day that the brace is worn. Patients have even been known to place objects into unworn braces of this type in order to fool the sensor. Coupled with the inconsistent patient compliance with brace usage, is a feeling by many physicians that braces, even if used properly, are not at all effective at curing scoliosis. These physicians may agree that bracing can possibly slow down or even temporarily stop curve (Cobb angle) progression, but they have noted that as soon as the treatment period ends and the brace is no longer worn, often the scoliosis rapidly progresses, to a Cobb angle even more severe than it was at the beginning of treatment. Some say the reason for the supposed ineffectiveness of the brace is that it works only on a portion of the torso, and not on the entire spine. Currently a prospective, randomized 500 patient clinical trial known as BrAIST (Bracing in Adolescent Idiopathic Scoliosis Trial) is enrolling patients, 50% of whom will be treated with the brace and 50% of who will simply be watched. The Cobb angle data will be measured continually up until skeletal maturity, or until a Cobb angle of 50° is reached, at which time the patient will likely undergo surgery. Many physicians feel that the BrAIST trial will show that braces are completely ineffective. If this is the case, the quandary about what to do with AIS patients who have a Cobb angle of between 20° and 40° will only become more pronounced. It should be noted that the “20° to 40°” patient population is as much as ten times larger than the “40° and greater” patient population.\nDistraction osteogenesis, also known as distraction callotasis and osteodistraction has been used successfully to lengthen long bones of the body. Typically, the bone, if not already fractured, is purposely fractured by means of a corticotomy, and the two segments of bone are gradually distracted apart, which allows new bone to form in the gap. If the distraction rate is too high, there is a risk of nonunion, if the rate is too low, there is a risk that the two segments will completely fuse to each other before the distraction period is complete. When the desired length of the bone is achieved using this process, the bone is allowed to consolidate. Distraction osteogenesis applications are mainly focused on the growth of the femur or tibia, but may also include the humerus, the jaw bone (micrognathia), or other bones. The reasons for lengthening or growing bones are multifold, the applications including, but not limited to: post osteosarcoma bone cancer; cosmetic lengthening (both legs-femur and/or tibia) in short stature or dwarfism/achondroplasia; lengthening of one limb to match the other (congenital, post-trauma, post-skeletal disorder, prosthetic knee joint), nonunions.\nDistraction osteogenesis using external fixators has been done for many years, but the external fixator can be unwieldy for the patient. It can also be painful, and the patient is subject to the risk of pin track infections, joint stiffness, loss of appetite, depression, cartilage damage and other side effects. Having the external fixator in place also delays the beginning of rehabilitation.\nIn response to the shortcomings of external fixator distraction, intramedullary distraction nails have been surgically implanted which are contained entirely within the bone. Some are automatically lengthened via repeated rotation of the patient's limb. This can sometimes be painful to the patient, and can often proceed in an uncontrolled fashion. This therefore makes it difficult to follow the strict daily or weekly lengthening regime that avoids nonunion (if too fast) or early consolidation (if too slow). Lower limb distraction rates are on the order of one mm per day. Other intramedullary nails have been developed which have an implanted motor and are remotely controlled by an antenna. These devices are therefore designed to be lengthened in a controlled manner, but due to their complexity, may not be manufacturable as an affordable product. Others have proposed intramedullary distractors containing and implanted magnet, which allows the distraction to be driven electromagnetically by an external stator. Because of the complexity and size of the external stator, this technology has not been reduced to a simple and cost-effective device that can be taken home, to allow patients to do daily lenthenings."}
{"text": "In the automotive industry, it is desirable to lock the steering shaft of a vehicle in place to prevent rotation of the steering shaft when the vehicle is not in use. Known systems utilize a keyed ignition system associated with the steering column. When the key is removed from the ignition, a steering shaft lock mechanism may be actuated to lock the steering shaft in place. Other vehicles utilize a keyless system or a system where the keyed ignition system is not associated with steering column. For these configurations, a separate actuator may be used to lock and unlock the steering shaft.\nIn one known steering shaft lock actuator, a locking member may be moved by the actuator from an unlocked to a locked position. When in the locked position, the locking member may engage a recess in the steering shaft to lock the steering shaft. The recess may be formed by adjacent teeth of a toothed ring that may be coupled to the steering shaft. An obstruction, e.g., a tooth of the toothed ring, may be encountered as the actuator attempts to drive the locking member to the locked position. The actuator motor may remain energized so that the motor drives the locking member to the locked position when the obstruction is removed. This can stress the motor and lead to early failure of the actuator. Actuator failure may be manifested in a failure of to lock the steering shaft when the vehicle is not in use and/or inadvertent locking of the steering shaft during normal driving conditions, which, of course, would pose a critical safety hazard.\nAccordingly, there is a need in the art for a steering shaft lock actuator configured to protect the actuator if an obstruction is encountered when the actuator attempts to drive the locking member toward a locked position and there is a need for a reliable steering shaft lock actuator to protect against inadvertent locking of the steering shaft during normal driving conditions.\nAlthough the following Detailed Description will proceed with reference being made to illustrative embodiments, many alternatives, modifications, and variations thereof will be apparent to those skilled in the art. Accordingly, it is intended that the claimed subject matter be viewed broadly."}
{"text": "Up-to-date motor vehicles, in particular passenger vehicles, are generally equipped with hydraulic or electrohydraulic servo steering systems, wherein a steering wheel is compulsively coupled mechanically with the steerable vehicle wheels. The servo assistance is devised such that actuators, e.g. hydraulic cylinders, are arranged in the mid-portion of the steering mechanism. A force generated by the actuators assists in the actuation of the steering mechanism in response to the turning of the steering wheel. This reduces the force the driver has to apply during the steering operation.\nSuperimposed steering systems are known in the art. They are characterized in that the steering angle input by the driver can be overridden in case of need by another steering angle (additional steering angle) by means of an actuator. Usually electric actuators are employed which act on an overriding drive and adjust the additional steering angle largely independently of the driver.\nThe additional steering angle is controlled by an electronic controller and is e.g. used to increase the stability and agility of the vehicle. According to a prior art control concept, as described in DE 197 51 125 A1, the steering components of the superimposed steering angle are produced irrespective of each other."}
{"text": "Engine oil tubes and fittings may be subjected to relatively high temperatures. Once subjected to excessive heating, oil may undergo coking Oil coking may cause solid oil deposits to form within oil tubes, causing undesirable effects such as blocked passageways and filters."}
{"text": "There is known a developing cartridge including a developing roller. The developing cartridge is detachably attachable to an image-forming apparatus.\nFor example, a conventional developing cartridge having an electrode for supplying power to a developing roller. The electrode has a shaft that can contact a power-supply section of an image-forming apparatus. This developing cartridge further includes a gear rotatable about the shaft and a protrusion provided on the gear. The protrusion can rotate together with the gear and is configured to contact a lever of the image-forming apparatus. The lever is moved when the protrusion contacts the lever. The image-forming apparatus can detect movement of the lever to thereby determine a specification of the developing cartridge attached to the image-forming apparatus."}
{"text": "The present invention relates to a device and a method for determining mixing ratios of flowing media, in particular for determining the mixing ratios of two gases by using two flow resistances with different characteristic curves.\nThe determination of the percentages of different components of fluids, for instance gas or liquid mixtures is necessary in many fields of engineering. For this purpose there exist a large variety of different methods using optical, acoustic, or electrochemical methods. Nonetheless, there is lacking a simple method for determining the mixing ratios of two fluids, in particular of two gases. Subject matter of the present invention is therefore a simple method for determining the relative portions of two components of a fluid mixture, in particular of a gas mixture, and a device suitable therefor.\nThus, the present invention teaches a device according to claim 1 and a method according to claim 6 carried out by means of said device. Advantageous embodiments of the invention are subject matter of subclaims referred thereto."}
{"text": "Voltage regulators are fundamental functions in powering most analog or digital functions in electronic systems. Operating at low currents, low power supply voltages, and over extended temperatures are key objectives in portable electronics systems, including in system-on-chips (SOC), especially in medical or emerging wireless and battery-less applications such as power-on-a-chip (PoC), Internet-of-things (IoT), or energy harvesting, to name a few.\nIn light of the potentially large markets for such emerging applications, unless solution costs are low, the markets may not realize their full potential reach. As such, low cost innovations in component designs are needed to enable such prospects of high volume markets to get off the ground and take off. In the last few decades, there has been a tremendous amount of annual manufacturing investment in semiconductors, some of which is sunk-cost. As a result, there is likely ample capacity in standard complementary metal-oxide-semiconductor (CMOS) manufacturing factories that may be readily and cheaply available for exploitation by new and promising high volume markets. The enablement of such new high volume markets would greatly benefit if their solutions require neither any custom devices nor any special processes, which could impose higher costs or risks on manufacturing yields. For design solutions in such emerging applications, to optimally exploit the existing digital CMOS capacity, manufacturing process node-portability should not be impeded. While the die-size of such solutions need to be kept small to keep the costs down, optimally no performance interference on the rest of the functions of the SOC should be imposed by innovative designs when serving such new applications requirements (low voltage and ultra-low current over maximal power supply and temperature spans).\nA survey of alternative voltage regulator topologies is presented in Fayomi, C J B. et al (2010) “Sub 1 V CMOS bandgap reference design techniques: a survey”, AICSP 62:141-157, DOI 10.1007/s10470-009-9352-4 discusses many techniques and their trade-offs for low power references and which is hereby incorporated by reference. Some of the trade-offs which are discussed in Fayomi et al. are briefly summarized as follows. Using Bipolar CMOS (BiCMOS) processes or depletion mode CMOS transistors frees up some operational supply headroom but adds to fabrication complexity and cost. Also, utilizing process deviations such as threshold differentials between positive channel metal-oxide semiconductors (PMOS) and negative channel metal-oxide semiconductors (NMOS), multiple depletion mode transistors, and differential thresholds via channel length variations. Alternatively, using custom devices such as a dynamic threshold metal-oxide semiconductor transistor (DTMOST), bulk driven CMOS, or lateral positive negative positive transistors (PNPs) can provide some design flexibility in extending digital power supply voltage (VDD) and bandgap voltage (VBG) spans, but at a price. Such deviations impose additional manufacturing costs, require special device modeling, or may demand post production control monitoring for the entire SOC that contains the reference. There are other design techniques, such as switch capacitors, that can provide extra operational headroom without the need for process or device variations. Switching techniques can help make low-noise references with stable temperature coefficients (TC) but may add cost due to capacitor extra mask or large capacitors for low droop rates. Chopping bandgap topologies can also cause high transient current consumption or may interfere with other substrate-noise sensitive analog functions within the SOC. Moreover, switch capacitor topologies can increase latch up risk and may require on chip charge-pumps in sub-1V environments. Banking only on CMOS transistors in subthreshold (which emulate pseudo-bipolar junction transistor (BJT) like behavior) to generate a VBG helps with low-power and sub-1V objectives, but such designs generally exhibit wider VBG variations, and likely require extra trimming and testing cost in post production. Forward biasing PMOS or NMOS body-source terminals can help increase bandgap amplifier headroom, but requires proper modeling at different process nodes. Additionally, such a design approach generally entails increased parasitic current and leakage risk at elevated temperatures, which limits the reference high temperature operations. The majority of reference topologies that impose on the SOC by requiring special processes or rely on non-standard use of transistors have not made it to high-volume applications. This is not just because of their added manufacturing cost, special device characterization, or special device modeling requirements. As noted earlier, variations on process, and deviation on devices for one function of an SOC (such as a reference) increases costs and manufacturing risk on the entire SOC and handicaps the SOC's process node portability, which ties the manufacturer's hands."}
{"text": "1. Field of the Invention\nThe present invention relates to a device and method for processing multimedia broadcast-related events in a portable terminal. More particularly, the present invention relates to a device and method for informing a user of a start of a frequently watched multimedia broadcast program using a character agent function.\n2. Description of the Related Art\nIn general, a portable terminal refers to a device which can be carried by a user, and its representative device is a mobile phone. Recently, the mobile phone has been developed so as to provide various additional functions as well as the original phone function. The mobile phone having the additional functions generally includes an LCD unit, and various display methods using character (avatar) image data have been developed in order to effectively display the additional functions and the state of the mobile phone.\nA DMB (Digital Multimedia Broadcasting) phone, through which satellite/terrestrial broadcasting can be viewed, is an example of a mobile phone having such various additional functions. The DMB is a broadcasting service by which broadcasting contents are transmitted so that various multimedia broadcasting such as video, audio, and data can be heard or viewed through multi-channels using a portable or vehicular receiver, to which a reception antenna is mounted, outside a house or while moving. When watching the DMB through a mobile phone, a user can select a desired broadcasting channel while watching various broadcasting channel guides displayed on the screen of the mobile phone, and can watch a program selected by the user by displaying a program of the selected channel on a display unit.\nA portable terminal including the multimedia receiving function receives CDM channel information related to the DMB service from a satellite. The received information includes EPG (Electronic Program Guide) information. Since the EPG information is the information related to broadcasting, and the table information of the EPG is continuously retransmitted by a predetermined interval, a user can receive the EPG information whenever the user uses the DMB service.\nFIG. 1 is an exemplary view for illustrating a construction of an EPG. Referring to FIG. 1, the EPG comprises a BIT (Broadcaster Information Table) 101, an SDT (Service Description Table) 103, an EIT (Event Information Table) 105, and a TOT (Time Offset Table) 107. Here, the BIT 101 is used to transmit broadcasting information. Here, notices, the names of broadcasting companies, and the current channel lists are transmitted. Broadcast IDs are used to discriminate between different BITs 101. Further, the SDT 103 is used to transmit information of the channels. The information on channels comprises the names of the channels, URL information on the channels, whether the channels are free or require a fee, and the types of the channels (whether the channels are video channels or audio channels, or whether the channels can be recorded). The SDT 103 is discriminated from another SDT using service IDs. Further, the EIT 105 represents genre information of the programs. Here, the name of the program, URL information of the program, a period of time from the start of the program to the completion of the program, and the viewer rating transmitted and the genre information of the corresponding program is transmitted. The EIT 105 uses event IDs to discriminate different events. The TOT 107 represents information on the present time. The TOT 107 is necessary to compare the time of the currently performed program with the present time or to display the present time. As noted above, a user can be informed of various DMB channel information through the EPG information.\nAn exemplary portable terminal including the multimedia reception function comprises a broadcast reservation function as well as a multimedia broadcast program watching function. The broadcast reservation function comprises a broadcast start alarm reservation function and a broadcast record reservation function. The broadcast start alarm reservation function is a function for informing a user of the start of a broadcast when a broadcast program starts, provided by registering the broadcast start alarm reservation of the broadcast program in a broadcast reservation menu, while the user is watching the EPG information through the screen of the terminal.\nHowever, according to the broadcast function, a user is required to set the reservation each time, and therefore, there are occasions in which the user does not reserve a favorite program either due to mistake, oversight or delay, and misses the start time of the program and can not watch the program. Further, the complexities involved with a user interface (UI) can make it inconvenient for a user to reserve a favorite program.\nAccordingly, a need exists for a system and method for easily and effectively registering a favorite multimedia broadcast program."}
{"text": "This invention relates generally to x-ray imaging processes and systems that use two or more x-ray spectra.\nConventional x-ray imaging uses a single kVp during a scan and represents an observed object in terms of its attenuation. However, this information is not sufficient to precisely characterize the observed object, since attenuation is energy dependent. In the diagnostic energy range, x-ray attenuation is the combination of two photon-matter interactions: the photoelectric effect and Compton scattering. These two interactions are energy-dependent and reflect the effective atomic number and electron density of the object. Accordingly, two measurements at distinct energy spectra can separate the attenuation information into these two basic components, which allows for better identification of the materials present in the object, contrast material specific images, or a weighted sum of the two interactions.\nLarger spectral separation of the low and high energy spectra is favorable to image variance deduction, therefore it contributes to reducing the total radiation dose of each scan while maintaining the same noise level of the output image."}
{"text": "The present invention relates to paper-manufacturing machines and to methods for operating such machines.\nIn particular, the present invention relates to twin-wire paper machines and methods for operating the same.\nMachines of this latter type conventionally include a pair of endless wires which respectively extend along closed loops, the wires having common run portions between which the web is compressed and formed while travelling from an inlet end to an outlet end of the common run portions.\nThe most commonly used type of paper-manufacturing machine is the fourdrinier machine which has a single wire guided so as to have an upper horizontal planar portion on which the web is formed. The forming of a continuous paper web in such a machine takes place on the substantially horizontal wire while the pulp slurry, or paper-fiber stock, is fed onto the wire from a headbox in the form of a dilute suspension in which water is present in a quantity of approximately 100-400 times the quantity of fibers. The water which is present in the fiber suspension is removed, or filtered off, through the wire in one direction, namely downwardly, partly due to gravity and partly by means of suction created by units which support the wire. The pulp stock is supplied to the fourdrinier from the headbox which has at a lower part adjacent the wire a narrow slice the length of which equals the width of the wire. From the slice the fiber suspension is discharged onto the wire in the form of a relatively coherent jet. The headbox is located at the starting end of the upper horizontal run of the wire and for the most part above the level of the wire where the web is formed.\nThe slice of the headbox is formed between upper and lower lips, and generally the upper lip is adjustable with respect to the lower lip so that it is possible in this way to regulate the angle at which the stock jet is directed toward the wire. The velocity of the stock jet with respect to the speed of the wire is also adjustable. Also a forming board which is customarily located beneath the wire between the breast roll and the first table roll is adjusted to determine the manner in which dewatering takes place during the initial stage of web formation, i.e. through a length of 50-150 cm subsequent to the slice.\nIn recent years there have been developments in connection with twin-wire paper machines which differ from fourdriniers in that the web is formed entirely between a pair of wires with dewatering taking place simultaneously in two opposed directions through both of the wires. This type of machine and principle of operation adds considerably to the dewatering rate. With such machines the area in which formation of the web takes place is generally curved with its direction deviating considerably from a horizontal plane. As a result of its efficient dewatering action, the twin-wire type of machine is considerably shorter than fourdrinier machines.\nWith known twin-wire machines the pulp stock is supplied by directing the stock jet from the headbox to the junction of the two wires at the inlet end of the common run portions thereof so that the stock jet enters at a converging entrance region where the wires converge toward each other to form a throat leading to the inlet end of the common run portions of the wires. The direction of the jet is maintained symmetrical with respect to the converging portions of the wires which travel toward the inlet end of the common run portions thereof, with the result that dewatering takes place rapidly in two opposed directions through both wires.\nWhile the above mode of feeding of the stock is without doubt efficient with respect to dewatering, nevertheless it involves, as compared with a fourdrinier, certain drawbacks as pointed out below.\nIn order to better understand the present invention and the problems solved thereby, certain theoretical factors in connection with formation of a paper web on a fourdrinier wire section are reviewed.\nThe paper produced on a paper machine always has a certain lack of homogeneity either with respect to properties which are visually observed or, for example, with respect to the strength characteristics of the paper. This lack of homogeneity partly results from the fact that the paper fibers are relatively long as compared with their thickness. As a result the fibers suspended in water become entangled, forming aggregates known as flocs which are visible as relatively dark spots in the finished paper when the latter is viewed against light. The presence of the flocs in paper is undesirable. Also as a result of the fact that the fibers are relatively long as compared to their thickness, the pulp suspension flows through the slice in such a way that the fibers tend to become oriented in the direction of flow and thus these fibers tend to remain oriented in the formed web longitudinally with respect to the long axis of the paper machine, which is to say the fibers tend to extend in the direction of web travel. This particular orientation of the fibers has a large influence on the properties of the strength of the paper as well as on the behavior of the paper with regard to water and water vapor.\nPaper which has its fibers oriented in a pronounced manner in the direction of web travel has a much higher tensile strength longitudinally, or in the machine direction, than in the cross-machine or transverse direction. In some cases such a high strength in the longitudinal direction is desirable, as when manufacturing spinning paper from which paper string is made. However, in most cases it is of advantage to maintain the tensile ratio of the paper, which is the ratio of the strengths thereof longitudinally and transversely, as low as possible.\nThe fiber orientation becomes clearly apparent from the behavior of the paper with respect to water and water vapor. When dry paper is transported, for example, into a humid storage area, it begins to curl as a result of the hygroscopic nature of the fibers as well as to the fact that the fibers swell to a greater extent transversely, as compared to any longitudinal swelling. Thus, where the paper has been manufactured in such a way that the fibers are primarily oriented in the direction of web travel, expansion due to humid conditions takes place in a much more pronounced manner transversely of the web, when considered with respect to the direction of machine travel, so that curling undesirably occurs.\nIt is therefore apparent that paper may be nonhomogeneous in two different ways. Flocs, which are created in the suspension and present in the paper formation, are immediately apparent to the eye. Anisotropy, or the difference in the characteristics of the paper longitudinally as compared with its characteristics transversely, is a phenomenon which may become visible under certain circumstances as when the paper curls in the manner described above. For the most part, however, this anisotropy is demonstrable by laboratory tests such as tests for determining strength properties. It is possible to influence both the presence of flocs as well as web formation and anisotropy by controlling the interaction between the headbox and the wire which receives the stock jet therefrom. The most important variables which should be controlled in this connection are:\n1. the relationship between the velocity of discharge of the stock jet and the speed of the wire, PA1 2. the angle with which the stock jet impinges on the wire, and PA1 3. the rate at which water is removed from the suspension which is deposited on the wire.\nThese factors affect not only the orientation of the fibers but also the web formation. In addition it is to be noted that if it is desired to reduce the presence of flocs and improve web formation, it is necessary to operate under conditions which will result in relatively pronounced fiber orientation. The paper-maker must choose operating conditions at the paper machine which will provide a compromise between those factors which will reduce fiber orientation and those factors which will improve web formation.\nBy way of example it should be noted that as a rule fiber orientation is least pronounced when the wire speed approximately equals the velocity with which the stock jet is discharged onto the wire. On the other hand, the best possible web formation requires the wire to travel at a speed somewhat higher than the jet velocity.\nIn addition, it is usually advisable, in order to prevent pronounced fiber orientation, to provide for removal of water from the deposited fiber suspension at a comparatively slow rate during the first few meters of web formation immediately subsequent to issue of the jet from the headbox. The dewatering process is influenced, for example, by regulating the angle at which the stock jet impinges on the wire. The smaller this angle the slower the passage of water through the wire. The magnitude of this angle can be changed most simply by adjusting the position with respect to each other of the lips of the headbox which define the slice thereof.\nHowever, slow dewatering impairs web formation in many cases.\nThus, there are no generally valid rules with respect to the influence of the above process variables on the characteristics of the paper. Optimization of paper machine operations depends, for example, on the fiber material which is utilized, the consistency of stock, and the machine speed, as well as, naturally, on the type and basis weight of the paper which is manufactured.\nThe achievement of the best results requires extensive experimentation and a paper machine which has the greatest possible versatility in the adjustment of the components thereof.\nWhile it is indeed possible with a fourdrinier to influence the fiber orientation during the initial 50-150 cm region of web formation on the wire beyond the slice of the headbox in the direction of web travel, such control of the fiber arrangement is not possible with twin-wire machines where it has been noted from observation and experience that when the pulp jet is directed into the throat to the inlet end of the common wire portions the arrangement and orientation of the fibers in the web which is produced is the same as in the jet itself. Further influence on fiber orientation is no longer possible with a twin-wire machine.\nOf course, it is recognized that many headbox designs presently used in twin wire machines deliver the pulp suspension into the gap between the wires in such condition that there is no particular fiber orientation in the jet and as a result the web which is formed has no pronounced fiber orientation. Nevertheless, as has been indicated above, there are certain types of paper which in view of their intended use should have a predetermined fiber orientation. For example orientation of the fibers longitudinally, in the direction of web travel, is highly favorable for certain intended uses of paper. In principle, orientation of the fibers may be achieved by discharging the stock from the slice at a speed smaller than the speed of travel of the wire on which the web is formed, in which case the combing effect of the wire which occurs at the start of the web formation results in fiber orientation. However with twin-wire machines no such combing effect can be produced because the arrangement of the fibers in the stock jet is \"frozen\" between the wires.\nFurthermore, with conventional twin-wire machines it is not possible to adjust to a desired extent the angle at which the stock jet impinges on the wire or wires, so that as a result of this latter deficiency one possible way of controlling web structure is lost."}
{"text": "1. Field of the Invention\nThe present invention relates to a vehicle surroundings monitoring apparatus for recognizing traveling circumstances in front of an own vehicle by stereoscopic cameras, monocular cameras, millimeter wave radars, and the like and for accurately estimating traveling paths of an own vehicle and a traveling control system incorporating such a vehicle surroundings monitoring apparatus.\n2. Discussion of Related Arts\nIn recent years, such a traveling control system as detecting traveling circumstances in front of an own vehicle by a camera and the like mounted on a vehicle, estimating traveling paths of the own vehicle from the traveling circumstances data, detecting a preceding vehicle traveling ahead of the own vehicle and making a follow-up control of the preceding vehicle or an intervehicle distance control between the own vehicle and the preceding vehicle, has been put into practical use.\nFor example, Japanese Patent Application Laid-open No. Toku-Kai-Hei 9-91598 discloses a traveling control system in which a traveling path of an own vehicle is estimated from traveling conditions such as yaw rate and other data and a nearest obstacle on the traveling path is detected as a preceding vehicle to be monitored. Further, in the traveling control system, when the preceding vehicle goes out of the traveling path of the own vehicle, the monitoring of the preceding vehicle is released.\nHowever, the prior technology in which a traveling path of an own vehicle (hereinafter referred to just as own traveling path) is estimated and a preceding vehicle is caught based on the own traveling path, has a defect that if the estimation of the own traveling path is inaccurate, the capture of the preceding vehicle itself loses reliability and as a result a desired traveling control can not be realized."}
{"text": "The invention relates to redundant control systems in which a plurality of redundantly arranged control units, such as computers, are monitored for fault conditions by circuitry that automatically responds to one or more failures by selecting a nonfaulty unit (or units) to control the system.\nSuch redundant control systems are used in a number of environments where failsafe operation is essential to the reliability and performance of the equipment. Digital flight controls in modern-day commercial and military aircraft are a prime example of one important application of such redundant systems. Other applications include traffic control, industrial plant or process control systems, spacecraft instrumentation, data processing systems, and telephone switching equipment.\nDigital flight control systems must, for apparent reasons, meet high reliability requirements. Normally, these requirements can be met only with redundant hardware, since it is virtually impossible to totally eliminate malfunctions of individual hardware components. In the management of redundant systems, there is a problem in distinguishing between operational and failed units and in selecting a proper configuration from the available, or remaining, healthy units. An automatic means and method for making this selection may be performed by logic circuitry that receives failure information signals from the redundant computers or other control units and makes the proper selection based on such information.\nIn developing the logic by which the selection circuitry chooses one control unit over another, a common design approach is to monitor the output states of the redundant control units and, by comparison voting of the outputs, determine by majority rule any faulty control units that should be excluded from participating in the control of the system. A variation of this logic scheme is to provide means associated with each control unit for cross-checking the health status of the remaining units and again voting the cross-checks so that when a majority of the control units determines that one or more minority control units have failed, those failed units will be excluded from system control. Systems based on such majority voting schemes work well as long as a majority of control units are available for voting down a minority of faulty units. When, however, the vote is even, such as where two remaining control units vote against each other as a result of cross-checks, then the system is at loggerheads and is unable to make a meaningful selection from the available fault information.\nOther proposed designs provide selection logic schemes that operate to select a certain control unit given a predetermined set of input fault conditions, even though another control unit is currently in control and has the same fault status. Such fixed, unflexible operating logic results in excessive switching between control units of equal probability of failure. Changing the unit in control, when unjustified by relative fault conditions, increases rather than decreases the likelihood of a malfunction in the overall system.\nStill another shortcoming of some existing selection logic circuits is their inability to recover from transient faults. Momentary fault conditions, including apparent (but not real) faults, may occur even though the monitored control unit is basically in good operating order. If the selection logic permanently excludes (deselects) a control unit because of a transient fault, the overall reliability of the redundant system rapidly deteriorates because the excluded control unit can no longer participate in a voting or other comparison scheme. On the other hand, if units are allowed to be brought back on line after they have recovered from a transient failure, then the overall system is better able to respond to future equipment failures."}
{"text": "For this surgical technique, only a limited amount of autologous spongiosa is available for filling the cavities of cage-like intervertebral or disk prosthesis and the spaces between individual implants and their surroundings. In the long term, the arthrodesis takes place not with the implant but between the bone and the bone replacement material. The individual implants therefore function only as place holders or spacers.\nThe intervertebral spaces, supplied with the known intervertebral implants, therefore frequently do not attain complete arthrodesis, that is, they end in a pseudoarthrosis. The situation is much the same also with cage-like intervertebral implants for the cervical spine, as well as for those, which were inserted through ventral entrances. Such intervertebral spaces are not stable mechanically, as would have been expected from a stiffening. The consequences then may be recurring pain with subsequent revision surgery.\nFor the implants and surgical techniques described above, the surgeon uses autologous bone material, which he obtains from the resected parts of the vertebral body or by means of an additional intervention in the crest of the ilium. Since dorsal accesses to the intervertebral disk space are very narrow, the applying of bone material is made difficult. The surgeon is unable to ensure that the whole of the intervertebral space is filled with autologous bone material. There is therefore the danger that empty spaces will result which, on the one hand, permits migration of the implant. On the other hand, the spaces, not filled with autologous bone material, are filled by a soft, fibrous tissue."}
{"text": "1. Field of the Invention\nThe present invention pertains to an improved marine propulsion system. In particular, the present invention pertains to improvements to a marine propulsion system of the type comprising a surface propeller, i.e., a propeller which is designed to be only partially submerged during operation so as to reduce the drag associated with the propeller. The propulsion system is a tractor-type system whereby the propeller is positioned ahead of the drive transmission housing so as to reduce the disturbances to the water upstream of the propeller, thereby improving the efficiency of the propeller. In addition, the present invention is designed with parallel drive and propeller shafts so that spur gears can be employed in the drive transmission, thereby eliminating unnecessary thrust bearings and improving the efficiency of the transmission gear train. As a further advantage of the present invention, the propulsion system is configured to enable easy retrofit of existing boat hulls and also the propeller is protected by the system.\n2. Description of the Related Art\nTractor marine propulsion systems have long been known in the art. For instance, Schmitt U.S. Pat. No. 1,605,376 discloses a prior art tractor drive propulsion system. This system comprises an engine mounted substantially forward of the boat transom and a long drive shaft extending aft from the engine through the bottom of the hull to a gearbox mounted below the transom. A propeller shaft extends forward from the gearbox at an angle relative to the drive shaft and a propeller is attached to the propeller shaft so as to be totally submerged beneath the hull. The drive shaft in this system is angled relative to the propeller shaft, thus the gearbox includes bevel gears. Due to the conical shape of the bevel gears, axial forces are produced which must be compensated for by thrust bearings or excessive wear and gear train inefficiencies occur. In addition, much of the Schmitt propulsion system is beneath the water. Thus, the brackets, shafts and gearbox increase the drag on the boat moving through the water over that of the hull alone. In addition, the shafts are exposed along much of their length, thereby suscepting them to damage from debris.\nBaldwin U.S. Pat. No. 3,105,455 also discloses a tractor drive propulsion system. The drive shaft and gearbox are mounted within the boat hull, thereby reducing drag and the potential for damage. In an alternate embodiment, Baldwin discloses a system having the motor substantially aft in the hull and the drive shaft extending forward to a universal Joint connected to the propeller shaft. The drive shaft and propeller shaft are angled relative to each other and the universal joint is needed to connect the two shafts. In both embodiments the propeller is completely submerged.\nLeavitt et al. U.S. Pat. No. 2,979,019 eliminated the angle between the tractor system drive shaft and propeller shaft by incorporating a belt and pulley transmission. However, the size of the housing below the waterline is enlarged to accommodate the pulley mounted to the propeller shaft, and thus the frontal area of the underwater structure is increased. This increase in frontal area increases the drag associated with the system and therefore decreases the overall efficiency. In addition, belt drive transmissions inherently cannot transmit as much power as the gearbox transmissions without slipping.\nAll the aforementioned prior art propulsion systems have relatively large frontal and surface areas below the waterline. Thus, the drag is increased and the propulsion system efficiency decreased. Arena U.S. Pat. No. 4,443,202 substantially reduces drag by employing a surface propeller propulsion system. The Arena system is of the pusher-type, i.e., the propeller is mounted aft of its drive shaft. The surface drive propeller is more efficient because less of the propeller drive system housing is submerged and therefore the drag associated with the system is decreased. However, a pusher-type system is often less efficient than a tractor system because of the disturbances created in the flow over the propeller by the drive transmission housing positioned forward of the propeller. In addition, the Arena system requires a significant cowling to protect personnel from the prop when entering and exiting the boat over the transom, as well as to shield over the rooster tail spray inherent in surface drive systems."}
{"text": "This specification relates to detecting light using a photodiode.\nLight propagates in free space or an optical medium is coupled to a photodiode that converts an optical signal to an electrical signal for processing."}
{"text": "1. Field of the Invention\nThe present invention relates to a board game and more particularly one in which a player establishes and builds a sports organization such as a baseball franchise club.\nEach player represents a baseball franchise and plays through a season in which the particular events occur that help or hinder his team and provide income or loss of income as a result of various events. These events are the ones that would normally occur to a team franchise throughout a season of play. Examples of the events are injuries to players, a player in a slump and sent to a minor league team, the purchase of a player from another team, contracting free agents, etc.\nThe board game simulates the real life events that occur to the team. A particularly novel feature of the present invention is the use of commercially available baseball cards.\n2. Description of the Prior Art\nA large number of board games exist representing sports events or business enterprises and examples of these are shown in U.S. Pat. Nos. 4,027,882 and 4,486,022. However, neither of these patents nor others in the prior art contemplate the use of commercially available baseball cards in which the game players represent baseball teams and the game contemplates playing through a season.\n3. Summary of the Invention\nThe game of the present invention is played on a board having a path along which playing pieces are moved. The path contains squares having indicia of the various events that can befall a baseball team throughout the season. Each player has a playing piece and in sequence rolls a die to advance his playing piece along the path on the board. As a playing piece comes to rest within a particular square, that player's team is subject to a designated event. It is understood that the events can either help or hurt a team and can provide income or expense to the club.\nThe purpose of the game is to enhance the quality of the team and receive income for the franchise. At the end of the game, the winner is determined by combining together the monetary income for the season plus points representing the quality of the team in terms of its players. The team quality is determined by the playing statistics of each player as represented on the commercial playing cards that each game player has at the end of the game.\nAccordingly it is a primary object of the present invention to provide a game board that stimulates the building up and playing of a baseball team during a season in which commercially available baseball cards are employed.\nAnother object of the present invention is to provide a baseball franchise game board using baseball cards representing real life baseball players in which each card contains the actual statistics of the player represented on the card.\nAnother object of the present invention is to provide a baseball franchise game board in which each game player represents a franchise baseball club and the object of the game is to build up and enhance the team by acquiring quality players represented by baseball cards of real life players and their respective baseball statistics."}
{"text": "Valves are used in myriad systems and environments, and may be operated in numerous and varied ways. Some valves are manually operated, others are operated via electrical, hydraulic, pneumatic, or various other types of actuators. In one particular implementation, a flow control valve is installed in an aircraft environmental control system. In such an implementation, the flow control valve is operated via a pneumatic actuation device that will open the flow control valve upon receipt of pressurized air upstream of the flow control valve. It may additionally be desirable, in such implementations, that the flow control valve include a shut-off mechanism that prevents downstream flow ducts from reaching or exceeding a particular pressure.\nThere is presently a desire for a control valve, which may be used in an aircraft environmental control system, that is normally-closed, is capable of being opened with pressure from either upstream or downstream of the control valve, is capable of being closed upon receipt of a suitable electrical signal, implements a backup means of closing if the electrical signal is not received, and is able to in the event that upstream pressure exceeds the electrical or backup closing setpoint. Unfortunately, there is no known valve that implements all of these functions.\nHence, there is a need for a flow control valve that implements the functionality described above. The present invention addresses at least this need."}
{"text": "1. Field of the Invention\nThe present invention relates to a sintered carbonitride alloy with Ti as main component and a cobalt binder phase, which has improved properties particularly when used as tool material for steel and cast iron cutting. More particularly, the present invention relates to a carbonitride-based alloy of specific composition and controlled relative saturation magnetization and coercivity for optimal combination of abrasive wear resistance, toughness and resistance to plastic deformation.\n2. Description of the Related Art\nTitanium-based carbonitride alloys, so called cermets, are widely used for metal cutting purposes. Compared to WC—Co based materials, cermets have excellent chemical stability when in contact with hot steel, even if it is uncoated, but have substantially lower toughness. This makes them most suited for finishing operations, which generally are characterized by limited mechanical loads on the cutting edge and a high surface finish requirement on the finished component.\nCermets comprise carbonitride hard constituents embedded in a metallic binder phase generally of Co and/or Ni. The hard constituent grains generally have a complex structure with a core, most often surrounded by one or more rims of other composition. In addition to Ti, group VIa elements, normally both Mo and W, are added to facilitate wetting between binder and hard constituents and to strengthen the binder phase by means of solution hardening. One or more of group IVa and/or Va elements, e.g. Zr, Hf, V, Nb and Ta, are also added in all commercial alloys available today. Cermets are produced using powder metallurgical methods. Powders forming binder phase and powders forming hard constituents are mixed, pressed and sintered.\nDuring recent years many attempts have been made to control the main properties of cermets in cutting tool applications, namely toughness, wear resistance and plastic deformation resistance. Much work has been done especially regarding the chemistry of the binder phase\nU.S. Pat. No. 6,344,170, U.S. Pat. No. 6,344,445 and U.S. Pat. No. 6,325,838 relate to a sintered body of a carbonitride alloy with titanium as main component with improved properties when used as cutting tool material. This has been achieved by combining a carbonitride based hard phase of specific chemical composition with an extremely solution hardened Co binder phase. By optimizing composition and sintering process in the Ti—Ta—W—C—N—Co system improved toughness and resistance to plastic deformation are accomplished. The two parameters that are used to optimize toughness and resistance to plastic deformation are the Ta- and Co-contents. The use of pure Co-based binder is a major advantage over mixed Co—Ni-based binders with respect to the toughness behavior due to the differences in solution hardening between Co and Ni.\nU.S. Pat. No. 7,332,122, and U.S. Pat. No. 7,157,044 are similar. They relate to a titanium based carbonitride alloy containing Ti, Nb, W, C, N and Co. By replacing Ta in U.S. Pat. No. 6,344,170 by Nb and carefully controlling the amount of undissolved Ti(C,N) cores a further optimization of technological properties has been achieved. More particularly, said patents relate to a carbonitride-based hard phase of specific composition, for which the amount of undissolved Ti(C,N) cores is optimized for maximal abrasive wear resistance, while the Co and Nb contents are simultaneously optimized to give the desired toughness and resistance to plastic deformation."}
{"text": "The invention relates generally to improvements in the lacing assembly of a shoe, and more particularly to an improved lacing assembly for an athletic shoe in which lower and upper vamp sections can be selectively adjusted and maintained to fit different portions of the foot independently while utilizing a single lace for and extending between the lower and upper vamp sections.\nIn applicant's prior U.S. Pat. No. 3,546,796, an athletic shoe is disclosed having separate lower and upper vamp sections. A pair of laces are provided, one lace for each vamp section. While each vamp section can be adjustably fitted and maintained by its own individual lace, it has been found to be disadvantageous to use a pair of separate laces for each shoe. For example, when untied, a pair of laces presents four loose ends that can become tangled and must be stored out into appropriate pairs before each vamp section can be adjustably fitted and tied. Moreover, when a pair of laces are tied, each lace provides a bow, thereby providing a pair of bows on each shoe that can become tangled and which present an unconventional and unacceptable appearance."}
{"text": "Computer networks are implemented using networking equipment such as, for example, optical networking ports, routers, and switches that are communicatively coupled with each other, often via physical media such as fiber optic cables. For example, a network may be implemented, in part, by coupling input/output ports of optical networking ports with router ports.\nAs the size of computer networks has expanded, the number of networking devices needed to implement those networks has increased and the interconnections between networking devices has become more complex. In the case of computer data centers which often comprise scores of computing resources, a high density of networking equipment and a complicated series of cable connections is often required to interconnect the numerous computing resources."}
{"text": "1. Field of the Invention\nThis invention relates to the field of acid-catalyzed organic reaction and particularly to methods of conducting acid-catalyzed reactions of organic compounds which reactions are promoted by strong acids. The invention also relates to novel acidic compositions useful in such reactions.\n2. Description of the Art\nThe ability of sulfuric acid to catalyze a variety of organic reactions is well known. It is also known that urea (a chalcogen compound useful in this invention) and sulfuric acid will combine to form adducts including the monourea-sulfuric acid adduct and the diurea-sulfuric acid adduct. For instance, D. F. du Toit, Verslag Akad. Wetenschappen, 22, 573-4 (abstracted in Chemical Abstracts, 8, 2346, 1914) disclosed that urea forms certain compounds with oxalic, acetic hydrochloric, nitric and sulfuric acids. L. H. Dalman, \"Ternary Systems of Urea and Acid. I. Urea, Nitric Acid and Water. II. Urea, Sulfuric Acid and Water. III. Urea, Oxalic Acid and Water\"; JACS, 56, 549-53 (1934), disclosed the phase relationships between the solid phase and saturated solutions containing urea and sulfuric acid at 10.degree. C. and 25.degree. C. The Sulfur Institute, Sulfur Institute Bulletin No. 10 (1964), \"Adding Plant Nutrient Sulfur to Fertilizer\", disclosed that urea reacts with sulfuric acid to form two complexes of \"urea sulfate\" which are useful fertilizers. Methods of manufacturing certain combinations of urea and sulfuric acid are disclosed by Verdegaal et al. in U.S. Pat. No. 4,310,343 and by Jones in U.S. Pat. No. 4,116,664.\nA wide variety of organic conversions are catalyzed by the proton-donating ability of strong acids. Such reactions have been extensively investigated and have been widely discussed in the literature. For instance, the Kirk-Othmer Encyclopedia of Chemical Technology, Third Edition, John Wiley and Sons, New York, 1980, discusses a variety of organic reactions that are catalyzed by strong acids including mineral acids, transition metal halides such as Friedel-Crafts catalysts, conjugate Friedel-Crafts catalysts, and others. Kirk-Othmer defines acid-catalyzed reactions as those in which a proton is transferred from the catalyst to the reactant which is thereby converted to an unstable state which immediately leads to the reaction under consideration. (Volume 5, page 33). While the proton donation mechanism of acid-catalyzed reactions referred to in Kirk-Othmer may or may not account for the reactions that take place in all acid-catalyzed reactions, it is known that strong acids promote numerous reactions including oxidative addition, reductive addition, esterification, transesterification, hydrogenation, isomerization (including racemization of optical isomers), hydrolysis and alcoholisis, alkylation, olefin polymerization, Friedel-Crafts reactions, demetalization of organics, and nitration reactions, among others. Strong acids known to be capable of promoting such acid-catalyzed organic reactions include sulfuric acid, nitric acid, hydrochloric acid, transition metal halides including the so-called Friedel-Crafts catalysts, for example, the halides of aluminum, gallium, boron, titanium, vanadium, tin and others, and conjugate Friedel-Crafts catalysts also known as Bronsted-Lewis superacid mixtures (Kirk-Othmer, V. 11, 295) such as mineral acid adducts of transition metal halides.\nAll of the known strong acid catalysts, and the methods involving their use for the promotion of acid-catalyzed organic reactions, suffer from one or more disadvantages. For instance, the strong mineral acids promote side reactions which form undesired by-products, destroy the organic feed material or product, and/or consume or deactivate the catalyst. Sulfuric acid is a strong sulfating, sulfonating, oxidizing, and dehydrating agent, and by virtue of those activities, it is consumed in most organic reactions by side reactions involving these mechanisms. Furthermore, the sulfonating and oxidizing activities of sulfuric acid result in the sulfonation and oxidation of organic feedstocks and/or products. Similar deficiencies exist with the other strong mineral acids such as hydrochloric and nitric acids. Hydrochloric acid chlorinates the reactants and thereby consumes the feed to produce unwanted chlorinated by-products. Nitric acid oxidizes and/or nitrates organic compounds. Hydrofluoric acid fluorinates organic reactants and products. The transition metal halides, including the Friedel-Crafts catalysts, are difficult to handle in that they must be isolated from water and reducing agents. Such catalysts also halogenate organic feedstocks and products.\nAccordingly, a need exists for improved methods of conducting acid-catalyzed organic reactions and for improved acid catalysts for use in such reactions which will promote the desired acid-catalyzed organic reaction yet reduce or eliminate the side reactions normally associated with acid-catalyzed organic reactions.\nIt is therefore a principal object of this invention to provide novel methods for the acid-catalyzed conversion of organic compounds.\nAnother object is the provision of novel methods for conducting acid-catalyzed reactions of organic compounds in the presence of compositions which comprise sulfuric acid.\nAnother object of this invention is the provision of novel acid catalysts comprising sulfuric acid which are effective for conducting acid-catalyzed organic reactions.\nAnother object of this invention is the provision of novel compositions which are useful for conducting acid-catalyzed organic reactions.\nAnother object of this invention is the provision of novel catalysts comprising sulfuric acid which have improved activity in the presence of lipophilic materials.\nYet another objective of this invention is the provision of novel methods for catalyzing organic reactions with sulfuric acid.\nAnother object is the provision of novel methods for the oxidative addition of organic compounds.\nYet another object is the provision of novel methods for reductive addition of organic compounds.\nAnother object is the provision of novel sulfuric acid-containing compositions useful for conducting organic reactions.\nAnother object is the provision of novel methods for the esterification and transesterification of organic compounds.\nYet another object of this invention is the provision of novel methods for hydrogenating organic compounds containing olefinic unsaturation.\nAnother object is the provision of novel methods for isomerizing organic compounds.\nYet another object is the provision of novel methods for the hydrolysis, alcoholisis, thiolosis, and amination of organic compounds.\nAnother object is the provision of novel methods for the alkylation of organic compounds.\nYet another object is the provision of novel methods for polymerizing olefinic compounds.\nYet another object is the provision of novel conjugate Friedel-Crafts catalysts.\nAnother object is the provision of novel Friedel-Crafts catalyzed organic reactions.\nYet another object of this invention is the provision of novel methods for demetalizing organic compounds.\nAnother object is the provision of novel methods for nitrating organic compounds.\nOther objects, aspects and advantages of this invention will be apparent to one skilled in the art in view of the following disclosure and the appended claims."}
{"text": "This invention relates to a sheet transport device for use in a printer, copying machine or like image forming apparatus.\nThere has been known a sheet transport device provided with feed rollers 3, 4 for dispensing sheets one by one from cassettes 1, 2 containing the sheets therein and a sheet separator 5 including a forward roller 18 which rotates in a direction along transport of a sheet and a retard roller 21 rotatable in a direction against the transport of the sheet as shown in FIG. 8. The retard roller 21 is mounted on a drive shaft which is biased toward the forward roller 18 by a spring or other biasing member, and accordingly is held in contact with the forward roller 18. Further, a torque limiter is mounted to the drive shaft of the retard roller 21 so as to control transmission of torque from a driving source to the drive shaft of the retard roller 21. The torque limiter is designed to cause the retard roller 21 to idly rotate upon receipt of the torque greater than a specified value.\nThe driving force of the forward roller 18 is set greater than that of the retard roller 21. In a normal sheet transport operation, the driving force of the transport roller 18 is transmitted to the retard roller 21 directly or through the sheet held between the rollers 18 and 21, activating the torque limiter to cause the retard roller 21 to idly rotate. Thereby, the retard roller 21 is caused to rotate in the sheet transport direction together with the forward roller 18 to transport the sheet.\nIn the case where a plurality of sheets are transported to a nip between the rollers 18 and 21 with one sheet adhered to another, i.e. in case of multiple feeding, the sheets come into sliding contact with one another, preventing transmission of the driving force of the forward roller 18 to the retard roller 21. This causes the retard roller 21 to rotate in a direction opposite to the rotating direction of the forward roller 18 and thereby the sheet in contact with the forward roller 18 and the one in contact with the retard roller 21 are subjected respectively to a transporting force and a retarding force working in opposing directions and separated from each other. As a result, transport of the sheet in contact of the retard roller 21 can be deterred.\nThe forward and retard rollers 18 and 21 are usually arranged as shown in FIGS. 9 and 10. With reference to FIG. 9, the forward roller 18 includes a long roller member while the retard roller 21 including a plurality of short roller members arranged at specified spacings along an axis. The number of roller members and the spacing between two adjacent roller members of the retard roller 21 are determined based on the width of copy sheets standardized for use in an image forming apparatus (hereafter referred to as standardized copy sheets). With this arrangement of the roller members of the retard roller 21, the following problems arise in the case where special sheets having the small width are to be transported. The special sheet may be material made of sheet, such as an envelope. Let it be assumed that special sheets A are being transported in line in a sheet transport device provided with a retard roller 21 including two roller members spaced along a drive shaft of the roller 21 as shown in FIG. 9. In this case, the roller members of the retard roller 21 unavoidably come into contact with the forward roller 18 at outer sides thereof. If only a portion of the retard roller 21 is in contact with the forward roller 18 when the sheet is separated, the driving force of the roller 18 is transmitted to the roller 21 and thereby the torque limiter causes the roller 21 to idly rotate. This will result in unsatisfactory sheet separation. In other words, since friction drag working on the contact portion of the rollers 18 and 21 is exceedingly large, the torque limiter disengages from the drive shaft of the roller 21 which in turn rotates together with the roller 18 in the sheet transport direction. Accordingly, in the case where two special sheets A are transported with the one sheet adhered to the other, i.e. in case of double feeding, the two special sheets A cannot be subjected to the transporting and retarding forces and therefore cannot be separated from each other properly.\nThe sheet transport device may be provided with a retard roller 21 having a single roller member as shown in FIG. 10 so as to properly separate special sheets A fed in line as described above. This sheet transport device suffers the disadvantage that, when the two special sheets A are transported side by side at one time, a center portion of the roller 21 comes into contact with the forward roller 18 and therefore the special sheets A cannot be separated properly.\nIn view of the above, a requirement is that the retard roller is not in contact with the forward roller when the special sheets having the small width such as envelopes are separated. Accordingly, it is preferable that the forward and retard rollers be away from each other in a normal state. However, since the standardized copy sheets which are most frequent in use are generally thin, it is preferable to keep the forward and retard rollers in contact with one another in the normal state in order to transport the standardized copy sheets reliably. On the contrary, the special sheets such as envelopes are generally thicker than the standardized copy sheets."}
{"text": "Iron-chromium alloys of substantial chromium content, such as stainless steel containing 12% or more chromium, are commonly made in a direct-arc electric arc furnace.\nSteel scrap is charged in the furnace hearth, the arc is struck and as soon as the arc is below the top scrap level, melt-down of the scrap is initiated under the maximum useful electric power. The scrap rapidly proceeds to melt, but as soon as melt-down proceeds to a degree where the arc energy is no longer adequately absorbed by unmelted scrap, power must be reduced to prevent the arc radiation from causing rapid deterioration of the furnace side wall and roof. After melt-down, ferro-chromium is added to the melt to obtain the required chromium content.\nIt is desirable to obtain a faster melt-down and eliminate the need for ferro-chromium which is a relatively expensive additive.\nThe cost of making such alloys would be reduced if the melt-down could proceed at a faster rate and the use of ferro-chromium eliminated."}
{"text": "Coated articles are known in the art for use in window applications such as insulating glass (IG) window units, vehicle windows, monolithic windows, and/or the like. In certain example instances, designers of coated articles often strive for a combination of high visible transmission, low emissivity (or low emittance), and/or low sheet resistance (Rs). High visible transmission may permit coated articles to be used in applications where these characteristics are desired such as in architectural or vehicle window applications, whereas low-emissivity (low-E), and low sheet resistance characteristics permit such coated articles to block significant amounts of IR radiation so as to reduce for example undesirable heating of vehicle or building interiors. Thus, typically, for coatings used on architectural glass to block significant amounts of IR radiation, high transmission in the visible spectrum is often desired. However, low transmittance and/or high reflectance in the IR and/or near IR part(s) of the spectrum is also desired to reduce for example undesirable heating of vehicle or building interiors.\nUnfortunately, low-E coatings often do not block significant amounts of ultraviolet (UV) radiation. In other words, low-E coatings typically provide only moderate or negligible UV protection, since the materials used in the layer stacks are transparent for short wavelengths (e.g., below 400 nm). In particular, materials used in such layer stacks such as tin oxide and titanium oxide cannot provide adequate UV protection given the small thicknesses of such materials required for low-E coatings. Thus, even with such coatings are provided on windows such as IG windows or vehicle windows, significant amounts of UV radiation makes its way through the window and into the building or vehicle. UV radiation tends to damage furniture and other elements inside of buildings or vehicles.\nMaterials such as vanadium oxide and cerium oxide absorb significant amounts of UV radiation. However, while such materials are characterized by a very steep onset of absorption for UV radiation, the onset of radiation absorption occurs in significant part in the visible part of the spectrum thereby leading to a significant distortion of colors when look through such a coating (e.g., a yellow shift). Accordingly, viewing characteristics tend to be degraded when layers of such materials are used.\nThere also exists a need in the art for improved chemical stability (chemical durability) and heat stability (stability upon heat treatment such as thermal tempering).\nIn view of the above, it will be appreciated that there exists a need in the art for a coated article including a low-E coating which is capable of blocking at some UV radiation in an efficient manner. Certain example embodiments of this invention relate to a coated article which permits significant UV absorption properties to be achieved.\nIn certain example embodiments of this invention, it has surprisingly been found that the provision of a layer consisting essentially of, or comprising, zirconium oxide (e.g., ZrO2), zirconium oxynitride, or zirconium silicon oxynitride (e.g., ZrSiOxNy) unexpectedly improves blocking (reflecting and/or absorption) of UV radiation in a manner which does not significantly degrade other optical properties of a coated article such as visible transmission and/or color. Surprisingly, when a layer comprising zirconium oxide or zirconium silicon oxynitride is provided as the uppermost or overcoat layer of the coated article (e.g., over a silicon nitride based layer), this results in improved chemical and heat stability in certain example embodiments.\nIn certain example embodiments of this invention, a layer of zirconium oxide or zirconium silicon oxynitride may be tuned in a manner so as to achieve a desired amount of UV blocking and/or absorption, as well as improved durability. It has been found that zirconium oxide or zirconium silicon oxynitride has optical constants (n and k) which allow adjustment of the onset of absorption by varying oxygen content of the layer for example. Moreover, it has been found that zirconium oxide, zirconium oxynitride, or zirconium silicon oxynitride has a refractive index (n) in a range which is very adaptable to low-E coatings, so that such layer(s) may be used in low-E coatings without significantly changing the visible appearance of the coated article or certain performance data. Thus, in certain example embodiments of this invention, the absorption edge of the curve defined by a layer of zirconium oxide or zirconium silicon oxynitride can be adjusted by changing the oxygen content thereof, which may be done for example by adjusting the amount of oxygen introduced into the sputtering chamber(s) during reactive sputter-deposition of the layer. In particular, for example, as oxygen content of the layer increases, the absorption edge of the curve defined by the layer of zirconium oxide or zirconium silicon oxynitride moves toward lower wavelengths away from certain visible wavelengths. Thus, in certain example embodiments, a balancing or tuning can be performed so as to achieve a desired balance between visible transmission and UV absorption.\nIn certain example embodiments of this invention, there is provided a coated article including a coating supported by a glass substrate, the coating comprising, in this order from the glass substrate outwardly: a first dielectric layer; a first contact layer; an infrared (IR) reflecting layer comprising silver located on the substrate over at least and contacting the first contact layer; a second contact layer comprising Ni and/or Cr located over and contacting the IR reflecting layer; a second dielectric layer comprising silicon nitride located over the second contact layer; and an overcoat dielectric layer comprises one or more of zirconium oxide, zirconium oxynitride, and/or zirconium silicon oxynitride located over and contacting the second dielectric layer comprising silicon nitride."}
{"text": "Response documents are commonly used in direct mail advertising to solicit responses from consumers. Examples of response documents include discount coupons, magazine subscription literature and product survey questionnaires. Typically, the addressee's name and address is recorded or coded on one or more response documents, the documents are mailed, and, ideally, the addressee then directly or indirectly returns the response document to the sender. Once returned, the response documents are used to determine which addressees responded, and, in some direct mail applications, how they responded.\nA typical direct mail application employs response documents in magazines to solicit magazine subscription renewals. An example of this type of application is described in Anderson et al. U.S. Pat. No. 3,819,173. In that patent, an addressee coded subscription card is inserted into an addressee-customized magazine while the magazine is being produced.\nIn a less common direct mail application, response documents in the form of discount coupons are coded with addressee information according to demographic information. This application entails selectively gathering different coupons for different addressees. The demographic information is used to determine which coupons would be most favorably responded to by each designated addressee. To effectively implement this technique, however, a substantial amount of manual assistance is required. For example, one such technique requires manually reviewing a list of coupons that are to be mailed to each addressee and then manually gathering the corresponding coupons. An addressee code is then stamped on each coupon before packaging and mailing the coupons to the addressee.\nRecently, there has been a need to mail large volumes of different response documents to different addressees. Additionally, for efficient processing, the different response documents that are sent to each designated addressee are required to be coded with information that identifies the particular addressee. Unfortunately, known direct mail techniques are not useful in meeting this recent need. For example, direct mail techniques that require a substantial amount of manual assistance, such as the technique described above, are intolerably burdensome. The volume of addressees and response documents renders such a technique inefficient and prohibitively expensive."}
{"text": "There currently exists a need for new techniques permitting the construction of power generating stations that are relatively small as compared with conventional reactor designs. For example, it is desirable that large-capacity power generating thermonuclear reactors be built in sizes of the order of 1000 times smaller than those possible with currently projected devices, to enable the utilization of these machines in spacecraft intended for deep space missions, or for utilization for marine propulsion. In these and other applications it may also be required to provide more direct conversion of thermonuclear energy to electrical energy than is available with conventional thermonuclear energy-converting apparatus.\nFor various additional reasons the present technology does not provide devices meeting the above criteria. Although heating plasma structures by magnetic mirror compression provides a basic solution, all previous attempts to do so have provided unsatisfactory results because the current methods involve heating a plasma ring that is moving with respect to the compression coils. This results in a very low coefficient of coupling and a very inefficient heating process. Furthermore, the rate of rise of the compression field in conventional devices must be very fast, requiring the use of expensive, complex, and bulky equipment.\nTherefore there currently exists a need for more efficient and less expensive means for compressing and heating significant amounts of plasma in thermonuclear devices. There is also a need for a multi-unit reactor assembly which is easier to maintain in operation or replace parts than in the currently available reactor designs.\nA typical background publication showing the present state of the art is U.S. Pat. No. 4,068,147, dated Jan. 10, 1978, to Daniel R. Wells."}
{"text": "1. Field of The Invention\nThis invention relates to catalyts and more particularly to improved methods of applying a catalytic material to an alumina catalyst.\n2. The Prior Art\nA calcined alumina coated metal mesh forms an effective support for many catalytic materials. Such catalysts can have many uses, such as an oxidizing catalyst useful for burning carbon monoxide or hydrocarbons in the exhaust gas streams of internal combustion engines, or as a self-regenerating filter for removing soot and other particulate matter from diesel engine exhaust gases. The formation of an alumina coated metal mesh and the addition of catalytic materials to the alumina coated metal mesh, if desired, is disclosed in many patents including co-assinged U.S. Pat. Nos. 3,227,659; 3,231,520; 3,240,698; 3,362,783; and 3,410,651. The alumina coated metal mesh can be formed by contacting a substrate, generally a metal mesh or other metallic article having an extended surface, with a coating solution which comprises an alkali metal aluminate as set forth in the prior art. After an adherent alumina film, generally of alumina trihydrate forms on the substrate, it is calcined to produce a porous alumina film. A catalytic material can then be added to the alumina film by one of several methods which have been disclosed in the prior art. The methods include applying the catalytic material to the calcined alumina film by vapor depostiton, contacting the alumina film with a solution of the catalytic material, or applying the catalytic material as a paste onto the surface of the alumina and then further treating the combination to produce the desired catalytic structure.\nGenerally, in the various processes of applying a catalytic material to the alumina film coated substrate, it is desired to uniformly coat the substrate and to produce a catalytic article of highest efficiency using the smallest possible quantity of catalytic material. This is because catalytic materials are generally quite expensive, and any saving in the quantity of catalytic material necessary to perform a desired function, results in a large decrease in the cost of the finished catalytic article. The various previously disclosed methods of applying a catalytic material to an alumina film coated substrate have had various degrees of success in producing the uniform distribution of catalytic material, and of producing an efficient catalyst."}
{"text": "1. Field of the invention\nThe present invention relates to a process for producing a polyhydric alcohol which is useful as a raw material for polyester resins, alkyd resins, polyurethane resins, polycarbonate resins, plasticizers, surfactants, lubricating oils, basis for cosmetics, reactive monomers and the like.\n2. Description of the Related Arts\nThere is known, as a process for producing a polyhydric alcohol, a two-stage reaction process in which an aliphatic aldehyde and formaldehyde are subjected to an aldol condensation reaction, and subsequently to a crossed Cannizzaro reaction each in the presence of a base catalyst, said aliphatic aldehyde being represented by the formula (i) ##STR1## wherein R.sup.1, R.sup.2 and R.sup.3 are each a hydrogen atom or a straight-chain or branched aliphatic hydrocarbon group having 1 to 22 carbon atoms with the proviso that at least one of them is a hydrogen atom.\nThe above-mentioned process is described in Japanese Patent Application Laid-Open No. 139141/1988 corresponding to U.S. Pat. No. 4,814,509, Japanese Patent Application Laid-Open No. 162538/1983 corresponding to U.S. Pat. No. 4,514,578 end the like, but each process described therein is based on the premise that the objective polyhydric alcohol is produced in combination with a formic acid salt.\nExamples of the base catalyst to be used in the above-mentioned process include a hydroxide of an alkali metal or an alkaline earth metal and a carbonate of any of the aforesaid metals, which are exemplified by sodium hydroxide, potassium hydroxide, calcium hydroxide, lithium hydroxide, sodium carbonate, potassium carbonate, calcium carbonate and lithium carbonate; and an amine compound, especially a tertially amine such as trimethylamine, triethylamine, diethylmethylamine, dimethylethylamine, diisopropylamine and tributylamine.\nAs a process for producing neopentyl glycol (hereinafter abbreviated to \"NPG\") which is one of polyhydric alcohols, there is described a process in which hydroxypivaldehyde (hereinafter abbreviated to \"HPA\") is synthesized at first by an aldol condensation reaction in the first reaction stage and then the resultant HPA is hydrogenated into NPG in Japanese Patent Application Publication No. 53421/1982 corresponding to EP 44412, Japanese Patent Application Laid-Open No. 299239/1989 corresponding to U.S. Pat. No. 4933493, and Japanese Patent Application Laid-Open No. 182442/1992 corresponding to U.S. Pat. No. 5,395,989, and the like patent.\nA process for producing trimethylol propane (hereinafter abbreviated to \"TMP\") which comprises effecting hydrogenation after the aldol condensation reaction in the same manner as above is described in Japanese Patent Application Laid-Open No. 92705/1978 corresponding to U.S. Pat. No. 412,290, Japanese Patent Application Laid-Open No. 287738/1988 corresponding to U.S. Pat. No. 5,149,861 and the like.\nThe production of a polyhydric alcohol in the presence of a base catalyst is accompanied with the by-production of a formic acid salt in an equimolar amount to the polyhydric alcohol as the main objective product. The formic acid salt as the by-product is utilized, for example, for the production of a hydrosulfite, etc., but is limited in its demand. Thus it follows that most of the by-produced formic acid salt is inevitably subjected to waste disposal, which incurs great expense, since it is accompanied by environmental pollution.\nOn the other hand, the process for producing NPG which comprises an aldol condensation reaction and successive hydrogenation requires a high pressure of about 50 to 150 kg/cm.sup.2 in the hydrogenation reaction, thereby complicating the production equipment.\nMoreover in the case of producing TMP by a process similar to the above, the objective TMP is obtained only at a low yield and less production rate per unit time; besides there has never been found an example demonstrating an industrial production of TMP by the foregoing process."}
{"text": "The invention relates in general to the field of computer-implemented methods and systems for managing datasets (e.g., files) in a storage system. In particular, it is directed to methods and systems for managing datasets across storage tiers of a storage system, based on their relevance.\nMulti-tiered storage systems may comprise several tiers of storage. Such systems typically assign different categories of data to various types of storage media, in order to reduce the global storage cost, while maintaining performance. A tiered storage system usually relies on policies that assign most frequently accessed data to high-performance storage tiers, whereas rarely accessed data are stored on low-performance (cheaper, and/or slower) storage tiers."}
{"text": "The present invention relates to a catalyst for the synthesis of chlorine dioxide. The catalyst is preferably a Pd/Al.sub.2 O.sub.3 or palladium and another platinum group metal (e.g., Pd+Pt/Al.sub.2 O.sub.3) or palladium and a Group IB metal (e.g., Pd+Au/Al.sub.2 O.sub.3) catalyst containing La.sub.2 O.sub.3 and Nd.sub.2 O.sub.3 which has increased activity and a slower rate of deactivation.\nChlorine dioxide is known to act as a disinfecting or sterilizing agent for solutions and devices (e.g., contact lenses). Chlorine dioxide is generally produced from sodium chlorite solutions by catalysts (e.g., catalysts containing noble metals, described in U.S. Pat. No. 5,008,096). However, known catalysts have the disadvantage of becoming greatly deactivated within a matter of days."}
{"text": "One prior art method for producing workpieces is known from DE 10 2005 057 945 A1. In such methods, rotationally symmetrical preforms are processed by pressure rollers to produce various products, such as, for example, support rolls, automotive parts and gas bottles. In view of the stress such products are exposed to, it is desirable to stiffen certain areas as early as possible when being formed. In utilization such products, in particular central sections thereof, are subjected to a stronger stress, so material reinforcement is preferable in this area. Designs with concave or conical inner surfaces are particularly preferred, however, they cannot be produced by known flow-forming methods.\nAccordingly, one object of the invention therefore is to provide a method of the above-mentioned character which enables the production of workpieces, from a perform, having a concave or conical inner surfaces."}
{"text": "The present invention relates to headlights of vehicles which are intended for use in difficult terrain. Such vehicles are mostly those for military use, such as tanks and halftracks which travel at considerable speed over difficult and irregular terrain.\nHeadlights of such vehicles are of rather complicated construction and comprise usually a metal housing the front of which is formed by a transparent lens-like wall seated in especially shaped, fluid and dust tight packings which latter are seated in appropriate seat grooves in the said casing. Furthermore, in view of the shocks and vibrations to which the headlights on such vehicles are exposed, special shock absorbing means have to be provided, being built into the lamp casing."}
{"text": "New wireless communication subscribers can register and activate their subscription, using their cellular phones in several ways.\nOne manner of activating a subscription is by activating a UICC such as a subscriber identification module (SIM) based mobile device in a wireless network. This activation consists in pre-programming the UICC with the mobile device with temporary activation identifiers, such as international mobile station identity (IMSI) and/or a mobile identification number (MIN), and a temporary electronic serial number which is used to identify the mobile device during registration and activation. The network identifies the temporary electronic serial number associated with the SIM vendor and invokes an over-the-air activation procedure for the SIM UICC mobile device. Once the activation process is completed, the temporary activation identifiers are overwritten.\nSuch activations, usually encountered on registration/activation processes are typically comprised in “second generation” (2G) or “third generation” (3G) of wireless services. Such activations typically rely on short message service (SMS) or push services, such as session initialisation protocol (SIP).\nLTE/4G mobile networks may comprise two layers: an IP packed-switched layer, on top of which is deployed an IP multimedia service (IMS) layer.\nThe IP packet-switched layer natively provides a TCP/IP connectivity between elements of the network. In particular, it natively enables TCP/IP communication between a UICC and a server.\nIMS provides a set of services as for example the voice over IP service, on top of the former IP packet-switched layer.\nThe IP Multimedia Subsystem (IMS) is a standardised networking architecture for the third generation and further generation mobile networks, which provides the users with mobile and fixed multimedia services. The IMS runs over the standard Internet Protocol (IP), using a Voice-over-IP (VoIP) implementation based on a 3GPP standardised implementation of Session Initiation Protocol (SIP). SIP, in turn, is a protocol developed for initiating, modifying, and terminating an interactive user session that involves multimedia elements such as video, voice, instant messaging, online games, and virtual reality.\nThis IMS set of services is the one that should be configured in order for the LTE subscription to be activated. However, with this method, a user still needs a push technology such as SIP or SMS for the OTA activation, and still needs to go in a point of sale (PoS) to activate the SIM card."}
{"text": "1. Field of the Invention\nThe present invention relates to an improvement of a solid state image pickup apparatus having a photoconductive film as a photoelectric converting section formed on a CCD scan substrate.\n2. Description of the Related Art\nThe essential section of a monolithic type CCD image pickup apparatus generally has the structure as shown in FIG. 1. Numeral 3 is an n-type buried channel on which transfer gates 6 and 7 are formed respectively through gate insulation films 5.sub.1 and 5.sub.2, thus constituting a vertical CCD. A photodiode 23 serving as a photoelectric converting section for the individual pixels is disposed adjacent to the vertical CCD. Numeral 2 is a p-type channel stopper. The pattern for one pixel portion of this CCD image pickup apparatus is as illustrated in FIGS. 2A-2C. FIG. 2A illustrates a signal charge transfer section 1-a between photodiode 23 and buried channel 3 as a shaded region, FIG. 2B buried channel 3 as a shaded region, and FIG. 2C an opening 18 of the photoelectric converting section constituted by photodiode 23, as a shaded region. As should be understood from these diagrams, since the monolithic type CCD image pickup apparatus has both of buried channel 3 and photodiode 23 formed on the same semiconductor substrate, the area of opening 18 is small. To reduce the lateral size of a pixel cell, therefore, the width of buried channel 3 decreases and so does the area of photodiode 23, as shown in FIGS. 3A-3C. This results in reduction in sensitivity of the image pickup apparatus and reduction in signal charge transfer capacity.\nIn contrast, a solid state image pickup apparatus having a photoconductive film formed on a CCD scan substrate has a signal charge transfer section separated from a photoelectric converting section, so that it is effective in realizing an image pickup apparatus with a high pixel density.\nFIG. 4 illustrates the structure of the essential section of such a laminated type CCD image pickup apparatus. N-type buried channel 3 is formed on a p-type Si substrate 1 and on the channel 3 are transfer gate electrodes 6 and 7 formed respectively through gate insulation films 5.sub.1 and 5.sub.2, thus constituting a vertical CCD. An n-type source layer 4 for reading out a signal charge of each pixel in buried channel 3 is disposed adjacent to the channel 3 of the vertical CCD. Since n-type source layer 4 is not used as photodiode 23 serving as the photoelectric converting section, it need not have a large area. Such a CCD scan substrate is covered with a first insulation film 8 in which a contact hole 9 is formed, and a relay electrode 11 that has a contact with each source layer 4 is formed in the hole 9. The resultant structure is covered with a second insulation film 12 in which another contact hole 13 is formed and a pixel electrode 14 is formed therein. Then, an undoped a-Si:H film 15 serving as a photoconductive film is disposed on the resultant structure, and a p-type a-SiC:H film 16 serving as a hole preventing layer is further disposed on the film 15, with a transparent electrode 17 being formed on the film 16.\nFIGS. 5A-5C illustrate patterns of a unit pixel region of this laminated type CCD image pickup apparatus in correspondence with FIGS. 2A-2C. The area of an opening 19 (see FIGS. 5C, 6C and 7C) determined by pixel electrode 14 is significantly larger than that of a monolithic type structure. With the illustrated structure, since photoconductive film 15 performs photoelectric conversion, that part which is other than the vertical CCD section for transferring a signal charge on the CCD scan substrate can have a small area. This can therefore permit the CCD image pickup apparatus to have such a layout that the area of buried channel 3 is made as large as possible by reducing the area of channel stopper 2, as shown in FIGS. 6A-6C. Accordingly, it is easy to increase the amount of signal charges transferred to 1.5 to 2 times as compared with a monolithic type CCD image pickup apparatus.\nIn such a laminated type CCD image pickup apparatus, however, reducing the lateral pixel size as shown in FIGS. 7A-7C considerably changes the width of buried channel 3, which adversely influences to signal charge transfer for the following detailed reason.\nFIG. 8 illustrates the results of measuring the channel potential (with zero gate voltage) with a change in channel width on a mask of a buried channel MOS transistor. FIG. 9 illustrates the structure of the MOS transistor, with the same numerals as used in FIGS. 1-4 denoting the corresponding sections. As should be obvious from FIG. 8, the channel potential of the buried channel becomes lower with a decrease in the channel width, and it can be hardly called a buried channel in a region with a channel width of less than 1 .mu.m. This is the so-called narrow channel effect. The measuring results show that, in the laminated CCD with a buried channel as shown in FIGS. 7A-7C, the channel potential at a region having a narrow channel width is low. With the transfer electrodes being arranged as shown in, for example, FIG. 10, the potential distribution at the C--C' cross section would be as illustrated in FIG. 11. More specifically, when the charges are transferred from a second layer transfer gate electrode 7-a to a region under a first layer transfer gate electrode 6, the distribution is as indicated by the one-dot chain line, and when charge transfer is from the electrode 6 to a region under a second transfer gate electrode 7-b, the distribution is as indicated by the broken line. In this state, a potential barrier leaves out some signal charge and the proper image reproduction cannot be provided. With the electrodes arranged as shown in FIG. 12, the potential distribution would be as illustrated in FIG. 13. In this case, incomplete charge transfer also occurs. To prevent such incomplete charge transfer, it is necessary to align the edge portions of first layer transfer gate electrode 6 with the boundaries at which the channel width varies, as shown in FIG. 14. This produces the potential distribution as shown in FIG. 15, which can prevent the aforementioned incomplete charge transfer. With such an electrode arrangement, however, the reverse charge transfer becomes impossible as per the case of two-phase driving structure in which a potential step is formed in one transfer gate electrode. For instance, an overflow drain is provided at the opposite edge of the vertical CCD to the horizontal CCD side, thus making the sweeping of the unnecessary charges impossible by the transfer function in the opposite direction to the transfer direction for reading signal charges.\nAs described above, according to the conventional laminated type CCD image pickup apparatuses, the narrow channel effect becomes prominent with a reduction in pixel cell size, and, particularly, a channel potential step is formed in one transfer gate electrode to thereby prevent bi-directional charge transfer."}
{"text": "1. Field of the Invention\nThe present invention relates to calix[4]arene dibenzo crown ether compounds, their preparation process and their use as extractants which can selectively separate cesium ions (Cs.sup.+). More particularly, the present invention relates to calix[4]arene dibenzo crown ethers able to selectively extract the cesium ions present in the state of traces in aqueous effluents from nuclear fuel cycle facilities, their preparation process and their use in extraction of cesium ions. Useful as ionophores for separating cesium ions from radioactive wastes, the compounds of the present invention are represented by the following structural formula 1: ##STR1## wherein R is a C.sub.1 -C.sub.10 normal alkyl, a phenyl or a methoxyphenyl, R' is hydrogen, p-tert-butyl or a C.sub.1 -C.sub.10 normal alkyl, R\" is hydrogen or a C.sub.1 -C.sub.10 normal alkyl, and n is an integer of 0-2.\n2. Description of the Prior Art\nAs industry has been highly advanced, environmental pollution is aggravated. In the present days, nuclear power is extensively used as an important energy source by virtue of its economical favorability in terms of, for example, investment cost in equipment, so more radioactive wastes are now generated mostly from atomic power plants. Therefore, the techniques for disposing radioactive wastes safely and in as little space as possible must be settled without delay.\nThe concentration technique by evaporation which has been employed for the disposal of radioactive waste liquids thus far, is very good in desalting, but produces a large amount of the wastes to be solidified because it treats all nuclides and salts together. In order to dispose of such wastes, of particularly interest are extractants which are capable of selectively separating heat-generating nuclides and long half-life nuclides. Active and extensive research has been and continues to be directed to the development of the extractants.\nA special control is required for cesium upon solidification after disposal of radioactive wastes because it is much longer in half-life (about 30 years) than other fission products, such as cobalt (5.72 years) and iodine (8 days), and generates heat. Thus, if the cesium ions contained in radioactive wastes are selectively separated, a great effect can be brought about on the disposal of radioactive wastes in terms of stability and energy efficiency.\nThe host-guest chemistry is a science field in which the interaction of the compounds, ions and/or molecules having intramolecular holes is researched. Since 1971 in which Helv. Chim. Acta, 54, 268, 1971, issued to professor Moff, describes that most of the ionophores present in nature are able to selectively transport cations through cell membranes, functioning as antibiotics, the chemistry has been of great interest and continues to be utilized in extensive research.\nCalixarenes, which have interesting structures, compose a branch of the host-guest chemistry. For the past twenty years, the research in the calixarene chemistry has been focused on the synthesis and structural specificity of calixarenes and their derivatives. Particularly, recent research has reported that the calixarene crown ether compounds based on calix[4]arenes are able to selectively extract the metals such as cesium and be applied for the solidification of the concentrated radioactive waste liquids, which could cause serious problems in the aspect of environmental pollution if they are not treated properly.\nThe word calixarene is the compound word of calix and arene, both derived from Greek, meaning a benzene ring-containing macrocyclic compound with a shape of a wine cup (Vogtel and Weber, Host Guest Complex Chemistry Macrocycles, 378, Springer-Verlag, 1985).\np-tert-Butyl calixarene, representative of calixarenes, was synthesized from the reaction of p-tert-butyl phenol and formaldehyde in the presence of sodium hydroxide, as shown in the following scheme 1 (Cornforth et al., J. Pharmacol., 73, 10, 1995): ##STR2##\nFrom the p-tert-butylcalix[4]arene 2, Gutsche et al., easily synthesized a calix[4]arene 3 by eliminating the p-tert-butyl group with the aid of an AlCl.sub.3 catalyst, as shown in the following scheme 2 (J. Org. Chem., 5795-5802, 50, 1985) ##STR3##\nD. N. Reinhoudt et al., succeeded in the synthesis of calix[4]arene-crown-6-ether and the isolation of its four stereoisomers. During a study on the use of these calix arene compounds as macrocyclic ligands to absorb alkali metals, they found that 1,3-alternate calix[4]arene-crown-6-ether is of high selectivity for cesium ions (J. Am. Chem. Soc., 117, 2767, 1995).\nDozol disclosed the synthesis of a biscrown calix[4]arene in PCT/FR93/01161 and a calix[4]arene crown ether in PCT/FR94/00432 and asserted that these macrocyclic ligands were superior in selectivity and efficiency for cesium ions to any other crown ether known. However, these crown ethers are problematic in practical use as an ion extractant because their binding capacity to ions or their ion-extracting performance is reduced."}
{"text": "1. Field of the Invention\nThe present invention relates to a receiver for receiving multi-channel broadcast and more particularly to a frequency converter which enables a VHF or UHF television receiver set to receive television broadcast waves using as their carriers superhigh frequency signals.\n2. Description of the Prior Art\nToday, the television broadcast uses electromagnetic waves belonging to VHF or UHF band, but the future trend in this field is toward the use of much higher frequencies, e.g. signal waves of SHF band.\nIn the practice of such a broadcasting service as using SHF band, it is of great economy if the currently used receiver set can be utilized in the reception of the broadcast. For this purpose, there is a need for a frequency converter which, if it is only attached to a conventional television receiver, makes the receiver adaptable for the new broadcast band. Above all, it is necessary to employ in the SHF broadcast the same television signals as the conventional ones and the same modulation system, i.e. vestigial-sideband amplitude modulation, as the one now in practice.\nThe frequency converter to enable the conventional TV receiver to receive the SHF broadcast, comprises, for example, an antenna for catching the broadcast waves in the air, a local oscillator for generating at a time a single oscillation frequency, a mixer for mixing the signals induced in the antenna and the output of the local oscillator, and an amplifier for amplifying the output, that is, frequency-converted signal, from the mixer. If the antenna catches three kinds of broadcast waves having their respective carrier frequencies f.sub.1, f.sub.2 and f.sub.3, the amplifier mentioned above will deliver the corresponding three kinds of signals having frequencies f'.sub.1, f'.sub.2 and f'.sub.3. Therefore, if the oscillation frequency of the local oscillator is so determined as to make the frequencies f'.sub.1, f'.sub.2 and f'.sub.3 belong to the conventional VHF or UHF band for TV broadcast, the SHF signal can be effectively received by simply connecting the converter with the conventional television receiver set. According to the convention of broadcast, the frequencies f'.sub.1, f'.sub.2 and f'.sub.3 are so determined as to correspond to the channels to which no broadcast waves are alloted in the district in consideration. The most important thing with such a frequency converter is the stability of the oscillation frequency of the local oscillator. The problem concerning the stability of the oscillation frequency will be described by the use of concrete numerical example.\nLet a case be considered in which an SHF broadcast channel having a video carrier frequency f.sub.1 of 12,000 MHz is received by converting the SHF channel to the fiftieth channel (video carrier frequency 693.25 MHz) of the present Japanese Standard TV Broadcast System. In this case, the local oscillation frequency f.sub.l, assumed to be set lower than the carrier frequency of the SHF channel, is such that EQU f.sub.l = 12,000 - 693.25 = 11,306.75 (MHz).\nThe video carrier frequency f.sub.1 of the broadcast wave can be considered almost fixed since the broadcast station is always monitoring and correcting the fluctuation in the frequency, so that the fluctuation in the frequency f'.sub.1 of the output of the converter depends mainly on the deviation of the local oscillation frequency f.sub.l, the variation +.DELTA.f of the frequency f.sub.l causing the variation -.DELTA.f of the frequency of the output of the converter.\nIn order to effectively receive television broadcast with conventional receiver sets, the video carrier frequency of the broadcast signal induced in the antenna must be confined within a range of at least the predetermined value .+-.0.1 MHz. Accordingly, the local oscillator of the frequency converter must have a frequency stability as high as 0.1/11,306.75, i.e. about 1 .times. 10.sup.-.sup.5. Such a high frequency stability can be attained by (1) using an oscillator having a high frequency stability, such as quartz or crystal controlled oscillator and (2) automatically controlling the local oscillation frequency in such a manner that the detected deviation of the output frequency is always rendered to zero.\nThe first method cannot be employed in the television sets for domestic use since the associated device is very complicated and expensive. Moreover, in the present level of technique, even a quartz controlled oscillator having the highest frequency stability attainable will not be able to maintain the above said stability in the circumstances varying in ambient temperature, power source voltage etc.\nThe second artifice, which has been employed in the domestic TV sets for the reception of VHF and UHF bands, cannot be applied, as it is, to the reception of the SHF television broadcast without encountering technical difficulties and above all without causing too much expense. This is described below.\nIn order to automatically control the local oscillation frequency, it is necessary to design the local oscillator of the frequency converter so that it may generate in a change-over fashion different frequencies in accordance with the number of the broadcast waves, and to incorporate in the frequency converter mentioned above a frequency discriminator which detects the frequency of the output of the mixer or the output of the amplifier and whose output is fed to the local oscillator to automatically control the oscillation frequency. In such a case, the carrier frequency in the output of the mixer remains to be a fixed value f.sub.o for any channel received since the local oscillation frequency is changed over in accordance with the frequencies of the received signals. Accordingly, the channels to be received are selected by changing over the oscillation frequencies of the local oscillator and the channel selecting device or channel selector of the conventional TV receiver set to which the frequency-converted signals are applied, is adjusted to the frequency f.sub.o.\nThe frequency converter described above has the following drawbacks.\n1. Since each of the plural TV broadcast channels will be separated from one another by about 12 MHz, as in the UHF band, the different frequencies generated by the local oscillator must differ from one another by about 12 MHz. As described before, however, the local frequency is to be around 11 GHz and it is technically very difficult to make a difference of 12 MHz with respect to such a high frequency.\n2. The condition necessary for the frequency stability proper to the local oscillator can be relaxed due to the use of the automatic frequency control, but the frequency excursion of the oscillator must be limited within a certain range so as to prevent the erroneous operations of the control system. Of all the erroneous operations, the one due to the influence of sound carrier takes place most usually. Namely, the frequency discriminator in its normal operation compares the video carrier frequency with the associated fixed value f.sub.o, but when the excursion of the local oscillation frequency exceeds a certain limit, the frequency discriminator will compare the value f.sub.o with the sound carrier frequency so that the control system operates erroneously. It is known from this that the allowable breadth of the excursion of the local oscillation frequency is about half (i.e. 2.3 MHz) the difference (4.5 MHz) between the video and the sound carrier frequencies. The breadth of 2.3 MHz is about 20 times the value 0.1 MHz obtained above in the absence of frequency control. The frequency stability in this case is 2 .times. 10.sup.-.sup.4 and this is a much relaxed condition in comparison with the previously obtained value of 1 .times. 10.sup.-.sup.5. In order to maintain the improved value under various conditions, however, the local oscillator must be provided with some auxiliary stabilizing means. Thus, under these circumstances, the local oscillator for generating not a single frequency but plural ones in change-over manner will add considerably to technical difficulties.\n3. As described above, the conventional TV receiver sets are utilized in the reception system using such a frequency converter as enlarged upon above. Since the receiver set has a channel selector, the separate provision of a channel selecting function in the converter will cause an operative complexity and above all be uneconomical."}
{"text": "Microelectronic integrated circuits (ICs) based on patterned semiconductor materials are continuing to evolve towards devices with an extremely high density of circuit elements per unit volume. The use of copper (Cu), with twice the conductivity of aluminum and three times the conductivity of tungsten, has been replaced as the interconnect material in advanced ICs manufacture using either single inlay or dual damascene processes. The conventional method of patterning trenches in a via-first dual damascene process involves filling via openings with an organic bottom anti-reflective coating (BARC), applying a photoresist layer on the via plug and then performing trench lithography. The via-first dual damascene process is considered superior in terms of reliable connection to the underlying metal, but has via proximity effect that causes difference in via CD (critical dimension).\nIn the semiconductor manufacturing art, the resolutions obtainable using optical lithography techniques is limited by what is known in the art as “proximity effects”, resulting in variations in feature of line width, alterations in pattern shapes, closed contacts or holes, and shortened or rounded lines among other distortions. Arising from proximity effects, the vias having difference critical dimensions, depending on their different pitches, usually occur and thereby cause defects. For example, a bird's beak is observed in a relative large via, and a blind via is observed in a relative small via. The most common technique to alleviate proximity effects is the use of optical proximity correction (OPC), which utilizes the addition of enlarged features on the mask pattern so as to anticipate and correct for the proximity effects. The OPC technique, however, is limited by the size of the added shapes. For compensating the via proximity effect, the use of OPC technique is difficult to make different-pitch vias obtain an identical via CD.\nTherefore, a method of alleviating proximity effects without the use of OPC technique is desired for the via-first dual damascene process. There remains, however, a need for a novel approach that affords an increased process window for the via patterning process."}
{"text": "A switching network includes a core module in a router and other data switching devices, and performs switching on a data packet or cell between multiple ports, that is, switches a data packet or cell, where the data packet or cell arrives at an input port, to a corresponding output port. A basic requirement on the switching network is that switched data packets and cells are required to keep an original order. As shown in FIG. 1, in a K-plane multi-level switching network in the prior art, normally a cell may arrive at a destination port through multiple paths, but it is very difficult for a time delay of each of the paths to be completely the same, which incurs an out-of-order phenomenon when cells that belong to a stream (from a same input port to a same output port) arrive at the destination port. The input port may be a source line card, and a component on the source line card, where the component on the source line card is directly related to the present invention, is an FIC (Fabric Interface Chip, switching network interface chip). The FIC performs processing on a data packet together with a switching network chip. An upstream FIC generally receives a slice of a variable-length data packet from a network side, and manages a virtual output queue (VOQ). The FIC may also include TM (Traffic Management, a traffic management function). An output port of the upstream FIC is a network card, and the network card includes mutually-independent switching chips or switching units. By taking a Clos interconnected multi-level switching network as an example, a switching chip may be formed by three levels, that is, switching units of three levels S1, S2, and S3. As shown in FIG. 2, a cell C1 departs from an FIC1, passes through an S1/3(2) and an S2(2), returns to an S1/3(2), and arrives at an FIC2; a cell C2 departs from the FIC1, passes through an S1/3(1), an S2(1), the S2(2), and the S1/3(2), and arrives at the FIC2; and a cell C3 departs from the FIC2, passes through the S1/3(1), the S2(2), the S2(1), and the S1/3(1), and arrives at the FIC1. If the cell C2 is sent first, and then the cell C1 and the cell C3 are sent at the same time, the cells C1 and C3 which are sent later may arrive at the S2(2) earlier than the cell C2 which is sent earlier, or the cells C1 and C3 which are sent later and the cell C2 which is sent earlier may arrive at the S2(2) at the same time, so that in this case an out-of-order problem occurs at the S2(2)."}
{"text": "Conventional methods of network security and content protection are often ineffective. Further limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "As an important chiral compound, (R)-3-hydroxyl-5-hexenoate has the following chemical formula (I), where R is alkyl or cycloalkyl having 1 to 8 carbon atoms, or mono- or poly-substituted aryl or aralkyl. The compound (I) has been widely used for synthesizing chemical and chemical intermediates. For example, Chen Fener et al. (WO2016074324) disclosed an efficient route for preparing, from derivatives of the compound (1), a key chemical intermediate (III) in synthesis of statin antilipemic agents.\n\nIt is indicated by document retrieval that chemical methods are main ways of synthesizing the compound (I) and structural analogues thereof at present. For example, U.S. Pat. No. 6,355,822 disclosed a method for stereoselectively preparing 3-hydroxyl-5-hexenoate (I) by reduction of a structural analogue 3-carbonyl-5-hexenoate (II) by using the resulting product of in-situ reaction of L-tartaric acid with sodium borohydride as chiral reductant. The reduction reaction needs to be conducted at a low temperature (−20° C. to −50° C.). Such reaction condition is strict, and the practical industrial application prospect is limited.\nEuropean Patent EP1176135 disclosed a preparation method for corresponding optically pure structural analogue 3-hydroxyl-5-hexenoate (I) by asymmetric catalyzation and hydrogenation of structural analogue 3-carbonyl-5-hexenoate (II) by using a complex of ruthenium and chiral phosphine ligand. Since this reaction needs to be conducted at a high temperature and under a high hydrogen pressure, it is inconvenient for industrial production.\nIt was yet reported that (R)-3-hydroxyl-5-hexenoate is selectively generated by biological catalysis. For example, referring to J. Chem. Soc. Perkin. Trans. 1, 1991, 133-140, it was reported that (R)-3-hydroxyl-5-hexenoate (1) was selectively generated by catalytic reduction of the compound (II) by baker's yeast, where the ee value was 78% (methyl ester) and 43% (ethyl ester). This method is less stereoselective and is not applicable to industrial production."}
{"text": "The present invention relates to a communication system, a storage device, and a control device for externally accessing a storage area storing file data, and particularly to a communication system, a storage device, and a control device for accessing file data handled in a file unit in a storage area of the storage device connected via a radio communication section.\nMore particularly, the present invention relates to a communication system, a storage device, and a control device that do not increase a process load on a storage device side at a time of access to file data, and particularly to a communication system, a storage device, and a control device that require addition of fewer functions necessary for access to file data on a storage device side.\nWith recent technological innovation, various information devices handling computer files of various media such as text, images, audio and the like have been developed and commercially available, the information devices including a personal computer, a PDA (Personal Digital Assistant), a digital camera, a portable media player and the like. Many of the information devices of this type have a connector or an interface for detachably mounting a cartridge type memory device such as a built-in memory, a hard disk, a memory stick, a memory card, a USB (Universal Serial Bus) memory or the like.\nFor example, a digital camera encodes a picked-up image into a predetermined file format such as JPEG (Joint Photographic Experts Group) or MPEG (Moving Picture Experts Group), and then stores the image on a memory card. In such a case, many images, which cannot be stored in a low-capacity built-in memory, can be stored. When the memory becomes full, a new capacity for storing picked-up images can be obtained by replacing the memory card.\nIn addition, files of images captured on the digital camera can be transferred to another information device such as a personal computer or the like to be subjected to image processing such as image quality adjustment, re-encoding and the like, and the images can be managed in the form of a database, an album or the like. Alternatively, an image file is transferred to a printer and printed out so that the image file can be appreciated as a photograph.\nIn a common method for transferring images and other data on the memory card between devices, the memory card is extracted and then loaded into a device to which to transfer the data. Alternatively, data transmission from the device in which the memory card is loaded to the device to which to transfer the data can be made via a wire or wireless transmission line using a wired interface or a wireless interface. The latter radio communication, in particular, eliminates a need for the operation of reattaching a connector and routing a cable each time data is transmitted, thus providing great convenience.\nFIG. 21 shows an example of image transmission by radio. In the figure, a digital camera is assumed as a mobile device. The digital camera stores picked-up images as image data in a built-in memory or an external memory card. In radio transmission, the digital camera reads desired image data from the built-in memory or the external memory card, and transfers the image data to image reproducing devices such as a personal computer, a television set, a printer and the like via a wireless interface. Of course, a wireless interface module is provided as an adapter in each of the receiving devices such as the personal computer, the television set, the printer and the like. After the transfer of the image data via a radio transmission line, the image data is displayed and stored by the personal computer, displayed on the screen of the television set, and printed out by the printer.\nAs described above, for radio transmission from the digital camera, a wireless interface module needs to be provided also on the transmission source device side. A device without a built-in wireless interface needs to be provided with a wireless interface by an external adapter. On the other hand, when a cartridge type storage device itself such as a memory stick or the like has a wireless interface module, the digital camera does not need to have a wireless interface module. In addition, without being attached to a data processing device proper such as the digital camera or the like, the memory stick can singly transmit data to a target device such as the personal computer, the printer or the like by radio.\nAn external memory medium having a radio communication function, for example, is proposed, for example, in Japanese Patent Laid-Open No. 2001-77878. This external memory medium includes a wire communication section for sending and receiving data to and from a host device into which the external memory medium is loaded via a physical connection section and a short-distance radio communication section for transmitting and receiving data to and from an external communication network via a short-distance radio communication network such as a Bluetooth communication network or the like. Hence, the host device having the external memory medium loaded therein can transmit and receive data to and from the communication network by driving the short-distance radio communication section within the memory medium.\nFIG. 22 is a diagram showing a common example of configuration of an electronic device to which a storage device such as a hard disk, a flash memory or the like is applied. The electronic device 4 shown in FIG. 22 includes a storage unit 43 and a function controlling unit 41.\nThe storage unit 43 includes a storage device 431 formed by a storage medium such as a magnetic disk, a memory element or the like and a storage device controlling functional unit 432 for making physical access to the storage medium.\nThe storage device 431 corresponds to a physical storage area itself. The storage device controlling functional unit 432 is a controller that performs control on the storage device 431 for faulty block management, error detection, error correction and the like. The storage unit 43 is controlled by the device controlling unit 41 via a predetermined electric interface 42 such as a processor bus.\nThe device controlling unit 41 is a controller for centralized control of operation of the device as a whole, which controller is provided by a microcontroller or a CPU (Central Processing Unit). The device controlling unit 41 has a file controlling unit 411 and a data controlling unit 412. The device controlling unit 41 reads and writes file data from and to the storage unit 43.\nThe data controlling unit 412 is a driver controlling unit that reads and writes data in a minimum physical access unit (generally 512 bytes) referred to as a “sector” or a “page” from and to the storage unit 43.\nThe file controlling unit 411 is generally referred to as a “file system,” and reads and writes file data, that is, data in a file unit from and to the storage unit 43 via the data controlling unit 412.\nIn related art, when the electronic device shares the data of the storage device with another device, a file system function and a data interface function (USB, IrDA, a wireless LAN or the like) are generally added to each device. FIG. 23 shows an example of configuration of the electronic device in such a case. The electronic device shown in FIG. 23 additionally includes an interface functional unit 44, and has an interface controlling unit 413 for driver control of the interface functional unit 44 in the device controlling unit 41.\nHowever, a device configuration as shown in FIG. 23 has problems of an increase in a process load on the device controlling unit 41 and an increase in power consumption due to the addition of the data interface functional unit. In addition, a load of driver development and the like for function introduction is imposed on a developer."}
{"text": "Acid sensing ion channels (ASICs) are proton-gated cationic channels expressed in the central and peripheral nervous systems. The ASIC family In humans includes ASIC1 , ASIC2, ASIC3 and ASIC4 subunits which arrange into homo- or heteromultimeric ion channels in neuronal membranes, ASIC1, ASIC2 and ASIC3 are significantly expressed in the small and medium nociceptive sensory neurons that are able to detect noxious chemical, thermal, and high threshold mechanical stimuli. ASIC2 and ASIC3 are also expressed in large neurons that mostly correspond to low threshold mechanoreceptors,\nASICs are permeable to Na+ and Ca2+ ions and are activated by external pH variations ranging from pH 6.8 to 4.0. ASICs are believed to play an important role in sensing local acidosis. Local tissue acidosis is a hallmark of pain and inflammation, and inflamed tissues frequently exhibit low pH (as low as ˜4.7). Blockage of ASICs has been proposed as a method for treating a variety of disorders and conditions. Blockage of ASIC1 in particular has been proposed as a means for treating conditions such as pain, neurodegenerative diseases, and psychiatric diseases.\nPharmacologic inhibitors of ASIC1 include the tarantula peptide Psalmotoxin-1 (PcTx1) which specifically inhibits ASIC1 homomers, and the small molecule, non-selective ASIC inhibitors amiloride and A-317567. The 40 amino acid peptide toxin APETx2, isolated from sea anemone, has been shown to inhibit ASIC3 homomers as well as ASIC3/1 and ASIC3/2 heteromers.\nCurrently, there are no known antibodies that specifically block ASICs. Thus, there is a need in the art for novel ASIC inhibitors, such as anti-ASIC1 antibodies, that can be used to treat diseases and conditions mediated by ASIC signaling."}
{"text": "The present invention relates to an oil separator for de-oiling crankcase ventilation gases of an internal combustion engine with the oil separator comprising a cyclone encompassing a gas inlet connected to the crankcase of the internal combustion engine, a gas outlet connected to the suction passage of the internal combustion engine, and an oil outlet connected to a crankcase sump of the internal combustion engine.\nAn oil separator of the kind mentioned above is known from DE 42 14 324 C2. In this oil separator a single cyclone is used which is dimensioned such that it may cope with the maximal amount of ventilation gases. An essential feature of this known oil separator is that a downward duct is connected to the oil outlet of the cyclone with the mouth thereof arranged below the oil level of the crankcase sump of the associated internal combustion engine, that a float valve which is normally open and arranged in a lower casing section at the side of the cyclone prevents a backwards flow out of the crankcase sump of the internal combustion engine, and that the complete oil separator as a construction unit is detachably connected to the downward pipe. These specific features of the oil separator have to take care that when the operating conditions of the associated internal combustion engine will certainly and strongly change, e.g., when full speed is achieved following an idle run, no oil from the crankcase sump of the internal combustion engine is transported into the interior of the cyclone because extremely large pressure differences will result from the change of the operating conditions.\nThe problem mentioned which is dealt with in the above document is solved by the oil separator described therein, however, this oil separator with the single cyclone has the disadvantage that no optimal oil separation is attained across the total range of operating conditions of the internal combustion engine which occur in practical operation. The occurring different operating conditions of the internal combustion engine on the one hand lead to different flow rates of ventilation gases out of the crankcase, and on the other hand to a different oil load of these ventilation gases. As a cyclone will operate only in a definite relative small range of operation conditions of the internal combustion engine, non-required oil volumes will occur in the gases out of the cyclone and fed to the suction air of the internal combustion engines during operation conditions outside of this optimal operating range.\nTherefore it is the object of the present invention to provide a generic oil separator avoiding the mentioned disadvantages and ensuring an optimal, i.e., complete or almost complete separation of oil from the crankcase ventilation gases in particular across a very large range of practical operation conditions of the internal combustion engine."}
{"text": "1. Technical Field\nThe present invention relates to an electromagnetic switch. In particular, the present invention relates to an electromagnetic switch provided with a single cylindrical frame and two electromagnets which are accommodated in the frame in parallel with each other in the axial direction of the solenoids.\n2. Related Art\nOne of conventional electromagnetic switches used for a starter of a vehicle is disclosed in Japanese Patent Laid-open Publication No. JP-A-2009-191843.\nThis electromagnetic switch has a first electromagnet for moving the pinion gear of the starter to a ring gear of an engine, and a second electromagnet for closing and opening a main switch. This switch additionally has a cylindrical frame in which the two electromagnets are set in parallel with each other in the axial direction of the switch, and a plastic cover which covers an open end of the frame and is fixed to the frame.\nFIG. 10 shows a schematic view of the electromagnetic switch viewed in the axial direction.\nTwo terminal bolts 110, 120 connected to a motor circuit are fixed to the plastic cover 100. A first exciting terminal 130 to be connected to a first coil of the first electromagnet is shown at the right side of the two bolts, and a second excitation terminal 140 to be connected to a second coil of the second electromagnet is shown at the left side of the two bolts in FIG. 10.\nThe first coil of the electromagnetic switch in JP-A-2009-191843 includes an attracting coil and a holding coil. In addition to the first and second excitation terminal 130, 140, the electromagnetic switch has a negative terminal 150 which is connected to a negative end of the attracting coil. The negative terminal 150 is projected from the plastic cover 100 to the axial direction, electrically connected to the terminal bolt 120 of the motor through an attachment clasp.\nAs shown in FIG. 10, the first and second excitation terminals 130, 140 face each other across the two terminal bolts 110, 120. Therefore it is necessary to wire up one of the excitation terminals 130, 140 before attaching the starter to the engine block. For example, when the engine block and the second excitation terminal 140 are arranged at the left side of the plastic cover 100 in FIG. 10, it is difficult to wire up the second excitation terminal 140 with the starter attached to the engine block. Accordingly, it is required to wire up the second excitation terminal 140 before attaching the starter to the engine block and to wire up the first excitation terminal 130 after attaching the starter to the engine block. This needs twice work of wiring before and after attaching the starter to the engine, and work efficiency is low.\nThe positional relationship of terminals to be wired up is different according to the engine model. It is necessary to change the position of the first and second excitation terminals 130, 140 according to the difference.\nHere a method of change the position of the first and second excitation terminals 130, 140 by rotating about an axis of the frame is discussed.\nThe terminal which is wired up after attaching the starter to the engine is defined as a basis position terminal (first excitation terminal 130 in the case described above). When the first and second excitation terminal 130, 140 are rotated about the axis with respect to the frame for changing the position of the basis position terminal from lower right to lower left, the positions of the first and second excitation terminal 130, 140 move from the position shown in FIG. 11 to the position shown in FIG. 12 with the positional relationship among the negative terminal, the first and second excitation terminal 130, 140 kept.\nIn the above described case, the negative terminal 150 is positioned in a contact space. The contact space is the space where a main switch is arranged. The main switch has a pair of fixed contacts which are connected to the two terminal bolts 110, 120, and a movable contact which closes and opens between the pair of fixed contacts. The electromagnetic switch in JP-A-2009-191843 is configured so that the first and second excitation terminals 130, 140 face each other across the contact space. In such an electromagnetic switch, rotating the first and second excitation terminals 130, 140 about the axis cannot change the position of the basis position terminal to the desired position often.\nIn the case of changing the position of the basis position according to the engine model and the like, for example as shown in FIG. 13 and FIG. 14, it is required to set the appropriate positional relationship between the first and second excitation terminals 130, 140 such that the other terminal wired up in advance (the second excitation terminal in FIG. 12) and the negative terminal 150 don't reach to the contact space. The change needs changing a position defined as a taking position for taking the first coil from a bobbin around which the first coil is wound. Therefore setting a plurality of variations according to different positions of terminals needs a plurality kind of bobbins according to the variations.\nThis cause the number of the component and the equipment used for setting the terminals in the electromagnetic switches according to a plurality of variations to increase. As a result, the production cost increases."}
{"text": "This invention relates to newly identified polynucleotides and polypeptides, and their production and uses, as well as their variants, agonists and antagonists, and their uses. In particular, in these and in other regards, the invention relates to novel polynucleotides and polypeptides of the tryptophanyl tRNA synthetase family, hereinafter referred to as xe2x80x9ctrpSxe2x80x9d.\nThe Streptococci make up a medically important genera of microbes known to cause several types of disease in humans, including, for example, otitis media, conjunctivitis, pneumonia, bacteremia, meningitis, sinusitis, pleural empyema and endocarditis, and most particularly meningitis, such as for example infection of cerebrospinal fluid. Since its isolation more than 100 years ago, Streptococcus pneumoniae has been one of the more intensively studied microbes. For example, much of our early understanding that DNA is, in fact, the genetic material was predicated on the work of Griffith and of Avery, Macleod and McCarty using this microbe. Despite the vast amount of research with S. pneumoniae, many questions concerning the virulence of this microbe remain. It is particularly preferred to employ Streptococcal genes and gene products as targets for the development of antibiotics.\nThe frequency of Streptococcus pneumoniae infections has risen dramatically in the past 20 years. This has been attributed to the emergence of multiply antibiotic resistant strains and an increasing population of people with weakened immune systems. It is no longer uncommon to isolate Streptococcus pneumoniae strains which are resistant to some or all of the standard antibiotics. This has created a demand for both new anti-microbial agents and diagnostic tests for this organism.\nt-RNA synthetases have a primary role in protein synthesis according to the following scheme:\nEnzyme+ATP+AAEnzyme.AA-AMP+PPi\nEnzyme.AA-AMP+t-RNAEnzyme+AMP+AA-t-RNA\nin which AA is an amino acid.\nInhibition of this process leads to a reduction in the levels of charged t-RNA and this triggers a cascade of responses known as the stringent response, the result of which is the induction of a state of dormancy in the organism. As such selective inhibitors of bacterial t-RNA synthetase have potential as antibacterial agents. One example of such is mupirocin which is a selective inhibitor of isoleucyl t-RNA synthetase. Other t-RNA synthetases are now being examined as possible anti-bacterial targets, this process being greatly assisted by the isolation of the synthetase.\nClearly, there is a need for factors, such as the novel compounds of the invention, that have a present benefit of being useful to screen compounds for antibiotic activity. Such factors are also useful to determine their role in pathogenesis of infection, dysfunction and disease. There is also a need for identification and characterization of such factors and their antagonists and agonists which can play a role in preventing, ameliorating or correcting infections, dysfunctions or diseases.\nThe polypeptides of the invention have amino acid sequence homology to a known Clostridium longisporum tryptophanyl tRNA synthetase protein.\nIt is an object of the invention to provide polypeptides that have been identified as novel trpS polypeptides by homology between the amino acid sequence set out in Table 1 [SEQ ID NO: 2] and a known amino acid sequence or sequences of other proteins such as Clostridium longisporum tryptophanyl tRNA synthetase protein.\nIt is a further object of the invention to provide polynucleotides that encode trpS polypeptides, particularly polynucleotides that encode the polypeptide herein designated trpS.\nIn a particularly preferred embodiment of the invention the polynucleotide comprises a region encoding trpS polypeptides comprising the sequence set out in Table 1 [SEQ ID NO:1] which includes a full length gene, or a variant thereof.\nIn another particularly preferred embodiment of the invention there is a novel trpS protein from Streptococcus pneumoniae comprising the amino acid sequence of Table 1 [SEQ ID NO:2], or a variant thereof.\nIn accordance with another aspect of the invention there is provided an isolated nucleic acid molecule encoding a mature polypeptide expressible by the Streptococcus pneumoniae 0100993 strain contained in the deposited strain.\nA further aspect of the invention there are provided isolated nucleic acid molecules encoding trpS, particularly Streptococcus pneumoniae trpS, including mRNAs, cDNAs, genomic DNAs. Further embodiments of the invention include biologically, diagnostically, prophylactically, clinically or therapeutically useful variants thereof, and compositions comprising the same.\nIn accordance with another aspect of the invention, there is provided the use of a polynucleotide of the invention for therapeutic or prophylactic purposes, in particular genetic immunization. Among the particularly preferred embodiments of the invention are naturally occurring allelic variants of trpS and polypeptides encoded thereby.\nAnother aspect of the invention there are provided novel polypeptides of Streptococcus pneumoniae referred to herein as trpS as well as biologically, diagnostically, prophylactically, clinically or therapeutically useful variants thereof, and compositions comprising the same.\nAmong the particularly preferred embodiments of the invention are variants of trpS polypeptide encoded by naturally occurring alleles of the trpS gene.\nIn a preferred embodiment of the invention there are provided methods for producing the aforementioned trpS polypeptides.\nIn accordance with yet another aspect of the invention, there are provided inhibitors to such polypeptides, useful as antibacterial agents, including, for example, antibodies.\nIn accordance with certain preferred embodiments of the invention, there are provided products, compositions and methods for assessing trpS expression, treating disease, for example, otitis media, conjunctivitis, pneumonia, bacteremia, meningitis, sinusitis, pleural empyema and endocarditis, and most particularly meningitis, such as for example infection of cerebrospinal fluid, assaying genetic variation, and administering a trpS polypeptide or polynucleotide to an organism to raise an immunological response against a bacteria, especially a Streptococcus pneumoniae bacteria.\nIn accordance with certain preferred embodiments of this and other aspects of the invention there are provided polynucleotides that hybridize to trpS polynucleotide sequences, particularly under stringent conditions.\nIn certain preferred embodiments of the invention there are provided antibodies against trpS polypeptides.\nIn other embodiments of the invention there are provided methods for identifying compounds which bind to or otherwise interact with and inhibit or activate an activity of a polypeptide or polynucleotide of the invention comprising: contacting a polypeptide or polynucleotide of the invention with a compound to be screened under conditions to permit binding to or other interaction between the compound and the polypeptide or polynucleotide to assess the binding to or other interaction with the compound, such binding or interaction being associated with a second component capable of providing a detectable signal in response to the binding or interaction of the polypeptide or polynucleotide with the compound; and determining whether the compound binds to or otherwise interacts with and activates or inhibits an activity of the polypeptide or polynucleotide by detecting the presence or absence of a signal generated from the binding or interaction of the compound with the polypeptide or polynucleotide.\nIn accordance with yet another aspect of the invention, there are provided trpS agonists and antagonists, preferably bacteriostatic or bacteriocidal agonists and antagonists.\nIn a further aspect of the invention there are provided compositions comprising a trpS polynucleotide or a trpS polypeptide for administration to a cell or to a multicellular organism.\nVarious changes and modifications within the spirit and scope of the disclosed invention will become readily apparent to those skilled in the art from reading the following descriptions and from reading the other parts of the present disclosure.\nThe following definitions are provided to facilitate understanding of certain terms used frequently herein.\nxe2x80x9cHost cellxe2x80x9d is a cell which has been transformed or transfected, or is capable of transformation or transfection by an exogenous polynucleotide sequence.\nxe2x80x9cIdentity,xe2x80x9d as known in the art, is a relationship between two or more polypeptide sequences or two or more polynucleotide sequences, as determined by comparing the sequences. In the art, xe2x80x9cidentityxe2x80x9d also means the degree of sequence relatedness between polypeptide or polynucleotide sequences, as the case may be, as determined by the match between strings of such sequences. xe2x80x9cIdentityxe2x80x9d and xe2x80x9csimilarityxe2x80x9d can be readily calculated by known methods, including but not limited to those described in (Computational Molecular Biology, Lesk, A. M., ed., Oxford University Press, New York, 1988; Biocomputing: Informatics and Genome Projects, Smith, D. W., ed., Academic Press, New York, 1993; Computer Analysis of Sequence Data, Part I, Griffin, A. M., and Griffin, H. G., eds., Humana Press, New Jersey, 1994; Sequence Analysis in Molecular Biology, von Heinje, G., Academic Press, 1987; and Sequence Analysis Primer, Gribskov, M. and Devereux, J., eds., M Stockton Press, New York, 1991; and Carillo, H., and Lipman, D., SIAM J. Applied Math., 48: 1073 (1988). Preferred methods to determine identity are designed to give the largest match between the sequences tested. Methods to determine identity and similarity are codified in publicly available computer programs. Preferred computer program methods to determine identity and similarity between two sequences include, but are not limited to, the GCG program package (Devereux, J., et al., Nucleic Acids Research 12(1): 387 (1984)), BLASTP, BLASTN, and FASTA (Atschul, S. F. et al., J. Molec. Biol. 215: 403-410 (1990). The BLAST X program is publicly available from NCBI and other sources (BLAST Manual, Altschul, S., et al., NCBI NLM NIH Bethesda, Md. 20894; Altschul, S., et al., J. Mol. Biol. 215: 403-410 (1990). As an illustration, by a polynucleotide having a nucleotide sequence having at least, for example, 95% xe2x80x9cidentityxe2x80x9d to a reference nucleotide sequence of SEQ ID NO: 1 it is intended that the nucleotide sequence of the polynucleotide is identical to the reference sequence except that the polynucleotide sequence may include up to five point mutations per each 100 nucleotides of the reference nucleotide sequence of SEQ ID NO: 1. In other words, to obtain a polynucleotide having a nucleotide sequence at least 95% identical to a reference nucleotide sequence, up to 5% of the nucleotides in the reference sequence may be deleted or substituted with another nucleotide, or a number of nucleotides up to 5% of the total nucleotides in the reference sequence may be inserted into the reference sequence. These mutations of the reference sequence may occur at the 5 or 3 terminal positions of the reference nucleotide sequence or anywhere between those terminal positions, interspersed either individually among nucleotides in the reference sequence or in one or more contiguous groups within the reference sequence. Analogously, by a polypeptide having an amino acid sequence having at least, for example, 95% identity to a reference amino acid sequence of SEQ ID NO:2 is intended that the amino acid sequence of the polypeptide is identical to the reference sequence except that the polypeptide sequence may include up to five amino acid alterations per each 100 amino acids of the reference amino acid of SEQ ID NO: 2. In other words, to obtain a polypeptide having an amino acid sequence at least 95% identical to a reference amino acid sequence, up to 5% of the amino acid residues in the reference sequence may be deleted or substituted with another amino acid, or a number of amino acids up to 5% of the total amino acid residues in the reference sequence may be inserted into the reference sequence. These alterations of the reference sequence may occur at the amino or carboxy terminal positions of the reference amino acid sequence or anywhere between those terminal positions, interspersed either individually among residues in the reference sequence or in one or more contiguous groups within the reference sequence.\nxe2x80x9cIsolatedxe2x80x9d means altered xe2x80x9cby the hand of manxe2x80x9d from its natural state, i.e., if it occurs in nature, it has been changed or removed from its original environment, or both. For example, a polynucleotide or a polypeptide naturally present in a living organism is not xe2x80x9cisolated,xe2x80x9d but the same polynucleotide or polypeptide separated from the coexisting materials of its natural state is xe2x80x9cisolatedxe2x80x9d, as the term is employed herein.\nxe2x80x9cPolynucleotide(s)xe2x80x9d generally refers to any polyribonucleotide or polydeoxribonucleotide, which may be unmodified RNA or DNA or modified RNA or DNA. xe2x80x9cPolynucleotide(s)xe2x80x9d include, without limitation, single- and double-stranded DNA, DNA that is a mixture of single- and double-stranded regions or single-, double- and triple-stranded regions, single- and double-stranded RNA, and RNA that is mixture of single- and double-stranded regions, hybrid molecules comprising DNA and RNA that may be single-stranded or, more typically, double-stranded, or triple-stranded regions, or a mixture of single- and double-stranded regions. In addition, xe2x80x9cpolynucleotidexe2x80x9d as used herein refers to triple-stranded regions comprising RNA or DNA or both RNA and DNA. The strands in such regions may be from the same molecule or from different molecules. The regions may include all of one or more of the molecules, but more typically involve only a region of some of the molecules. One of the molecules of a triple-helical region often is an oligonucleotide. As used herein, the term xe2x80x9cpolynucleotide(s)xe2x80x9d also includes DNAs or RNAs as described above that contain one or more modified bases. Thus, DNAs or RNAs with backbones modified for stability or for other reasons are xe2x80x9cpolynucleotide(s)xe2x80x9d as that term is intended herein. Moreover, DNAs or RNAs comprising unusual bases, such as inosine, or modified bases, such as tritylated bases, to name just two examples, are polynucleotides as the term is used herein. It will be appreciated that a great variety of modifications have been made to DNA and RNA that serve many useful purposes known to those of skill in the art. The term xe2x80x9cpolynucleotide(s)xe2x80x9d as it is employed herein embraces such chemically, enzymatically or metabolically modified forms of polynucleotides, as well as the chemical forms of DNA and RNA characteristic of viruses and cells, including, for example, simple and complex cells. xe2x80x9cPolynucleotide(s)xe2x80x9d also embraces short polynucleotides often referred to as oligonucleotide(s).\nxe2x80x9cPolypeptide(s)xe2x80x9d refers to any peptide or protein comprising two or more amino acids joined to each other by peptide bonds or modified peptide bonds. xe2x80x9cPolypeptide(s)xe2x80x9d refers to both short chains, commonly referred to as peptides, oligopeptides and oligomers and to longer chains generally referred to as proteins. Polypeptides may contain amino acids other than the 20 gene encoded amino acids. xe2x80x9cPolypeptide(s)xe2x80x9d include those modified either by natural processes, such as processing and other post-translational modifications, but also by chemical modification techniques. Such modifications are well described in basic texts and in more detailed monographs, as well as in a voluminous research literature, and they are well known to those of skill in the art. It will be appreciated that the same type of modification may be present in the same or varying degree at several sites in a given polypeptide. Also, a given polypeptide may contain many types of modifications. Modifications can occur anywhere in a polypeptide, including the peptide backbone, the amino acid side-chains, and the amino or carboxyl termini. Modifications include, for example, acetylation, acylation, ADP-ribosylation, amidation, covalent attachment of flavin, covalent attachment of a heme moiety, covalent attachment of a nucleotide or nucleotide derivative, covalent attachment of a lipid or lipid derivative, covalent attachment of phosphotidylinositol, cross-linking, cyclization, disulfide bond formation, demethylation, formation of covalent cross-links, formation of cysteine, formation of pyroglutamate, formylation, gamma carboxylation, glycosylation, GPI anchor formation, hydroxylation, iodination, methylation, myristoylation, oxidation, proteolytic processing, phosphorylation, prenylation, racemization, glycosylation, lipid attachment, sulfation, gamma-carboxylation of glutamic acid residues, hydroxylation and ADP-ribosylation, selenoylation, sulfation, transfer-RNA mediated addition of amino acids to proteins, such as arginylation, and ubiquitination. See, for instance, PROTEINSxe2x80x94STRUCTURE AND MOLECULAR PROPERTIES, 2nd Ed., T. E. Creighton, W. H. Freeman and Company, New York (1993) and Wold, F., Posttranslational Protein Modifications: Perspectives and Prospects, pgs. 1-12 in POSTTRANSLATIONAL COVALENT MODIFICATION OF PROTEINS, B. C. Johnson, Ed., Academic Press, New York (1983); Seifter et al., Meth. Enzymol. 182:626-646 (1990) and Rattan et al., Protein Synthesis: Posttranslational Modifications and Aging, Ann. N.Y. Acad. Sci. 663: 48-62 (1992). Polypeptides may be branched or cyclic, with or without branching. Cyclic, branched and branched circular polypeptides may result from post-translational natural processes and may be made by entirely synthetic methods, as well.\nxe2x80x9cVariant(s)xe2x80x9d as the term is used herein, is a polynucleotide or polypeptide that differs from a reference polynucleotide or polypeptide respectively, but retains essential properties. A typical variant of a polynucleotide differs in nucleotide sequence from another, reference polynucleotide. Changes in the nucleotide sequence of the variant may or may not alter the amino acid sequence of a polypeptide encoded by the reference polynucleotide. Nucleotide changes may result in amino acid substitutions, additions, deletions, fusions and truncations in the polypeptide encoded by the reference sequence, as discussed below. A typical variant of a polypeptide differs in amino acid sequence from another, reference polypeptide. Generally, differences are limited so that the sequences of the reference polypeptide and the variant are closely similar overall and, in many regions, identical. A variant and reference polypeptide may differ in amino acid sequence by one or more substitutions, additions, deletions in any combination. A substituted or inserted amino acid residue may or may not be one encoded by the genetic code. A variant of a polynucleotide or polypeptide may be a naturally occurring such as an allelic variant, or it may be a variant that is not known to occur naturally. Non-naturally occurring variants of polynucleotides and polypeptides may be made by mutagenesis techniques, by direct synthesis, and by other recombinant methods known to skilled artisans.\nThe invention relates to novel trpS polypeptides and polynucleotides as described in greater detail below. In particular, the invention relates to polypeptides and polynucleotides of a novel trpS of Streptococcus pneumoniae, which is related by amino acid sequence homology to Clostridium longisporum tryptophanyl tRNA synthetase polypeptide. The invention relates especially to trpS having the nucleotide and amino acid sequences set out in Table 1 [SEQ ID NO: 1] and Table 1 [SEQ ID NO: 2] respectively, and to the trpS nucleotide sequences of the DNA in the deposited strain and amino acid sequences encoded thereby.\nDeposited Materials\nA deposit containing a Streptococcus pneumoniae 0100993 strain has been deposited with the National Collections of Industrial and Marine Bacteria Ltd. (herein xe2x80x9cNCIMBxe2x80x9d), 23 St. Machar Drive, Aberdeen AB2 1RY, Scotland on Apr. 11, 1996 and assigned deposit number 40794. The deposit was described as Streptococcus pneumnoiae 0100993 on deposit. On Apr. 17, 1996 a Streptococcus pneumnoiae 0100993 DNA library in E. coli was similarly deposited with the NCIMB and assigned deposit number 40800. The Streptococcus pneumoniae strain deposit is referred to herein as xe2x80x9cthe deposited strainxe2x80x9d or as xe2x80x9cthe DNA of the deposited strain.xe2x80x9d\nThe deposited strain contains the full length trpS gene. The sequence of the polynucleotides contained in the deposited strain, as well as the amino acid sequence of the polypeptide encoded thereby, are controlling in the event of any conflict with any description of sequences herein.\nThe deposit of the deposited strain has been made under the terms of the Budapest Treaty on the International Recognition of the Deposit of Micro-organisms for Purposes of Patent Procedure. The strain will be irrevocably and without restriction or condition released to the public upon the issuance of a patent. The deposited strain is provided merely as convenience to those of skill in the art and is not an admission that a deposit is required for enablement, such as that required under 35 U.S.C. xc2xa7112.\nA license may be required to make, use or sell the deposited strain, and compounds derived therefrom, and no such license is hereby granted.\nPolypeptides\nThe polypeptides of the invention include the polypeptide of Table 1 [SEQ ID NO:2] (in particular the mature polypeptide) as well as polypeptides and fragments, particularly those which have the biological activity of trpS, and also those which have at least 70% identity to a polypeptide of Table 1 [SEQ ID NOS:2 and 4] or the relevant portion, preferably at least 80% identity to a polypeptide of Table 1 [SEQ ID NOS:2 and 4], and more preferably at least 90% similarity (more preferably at least 90% identity) to a polypeptide of Table 1 [SEQ ID NOS:2 and 4] and still more preferably at least 95% similarity (still more preferably at least 95% identity) to a polypeptide of Table 1 [SEQ ID NOS:2 and 4] and also include portions of such polypeptides with such portion of the polypeptide generally containing at least 30 amino acids and more preferably at least 50 amino acids.\nThe invention also includes polypeptides of the formula set forth in Table 1 (D) [SEQ ID NO:2] wherein, at the amino terminus, X is hydrogen, and at the carboxyl terminus, Y is hydrogen or a metal, R1 and R2 is any amino acid residue, and n is an integer between 1 and 1000. Any stretch of amino acid residues denoted by either R group, where R is greater than 1, may be either a heteropolymer or a homopolymer, preferably a heteropolymer.\nA fragment is a variant polypeptide having an amino acid sequence that entirely is the same as part but not all of the amino acid sequence of the aforementioned polypeptides. As with trpS polypeptides fragments may be xe2x80x9cfree-standing,xe2x80x9d or comprised within a larger polypeptide of which they form a part or region, most preferably as a single continuous region, a single larger polypeptide.\nPreferred fragments include, for example, truncation polypeptides having a portion of an amino acid sequence of Table 1 [SEQ ID NOS:2 and 4], or of variants thereof, such as a continuous series of residues that includes the amino terminus, or a continuous series of residues that includes the carboxyl terminus. Degradation forms of the polypeptides of the invention in a host cell, particularly a Streptococcus pneumoniae, are also preferred. Further preferred are fragments characterized by structural or functional attributes such as fragments that comprise alpha-helix and alpha-helix forming regions, beta-sheet and beta-sheet-forming regions, turn and turn-forming regions, coil and coil-forming regions, hydrophilic regions, hydrophobic regions, alpha amphipathic regions, beta amphipathic regions, flexible regions, surface-forming regions, substrate binding region, and high antigenic index regions.\nAlso preferred are biologically active fragments which are those fragments that mediate activities of trpS, including those with a similar activity or an improved activity, or with a decreased undesirable activity. Also included are those fragments that are antigenic or immunogenic in an animal, especially in a human. Particularly preferred are fragments comprising receptors or domains of enzymes that confer a function essential for viability of Streptococcus pneumoniae or the ability to initiate, or maintain cause disease in an individual, particularly a human.\nVariants that are fragments of the polypeptides of the invention may be employed for producing the corresponding full-length polypeptide by peptide synthesis; therefore, these variants may be employed as intermediates for producing the full-length polypeptides of the invention.\nPolynucleotides\nAnother aspect of the invention relates to isolated polynucleotides, including the full length gene, that encode the trpS polypeptide having a deduced amino acid sequence of Table 1 [SEQ ID NOS:2 and 4] and polynucleotides closely related thereto and variants thereof.\nUsing the information provided herein, such as a polynucleotide sequence set out in Table 1 [SEQ ID NOS:1 and 3], a polynucleotide of the invention encoding trpS polypeptide may be obtained using standard cloning and screening methods, such as those for cloning and sequencing chromosomal DNA fragments from bacteria using Streptococcus pneumoniae 0100993 cells as starting material, followed by obtaining a full length clone. For example, to obtain a polynucleotide sequence of the invention, such as a sequence given in Table 1 [SEQ ID NOS:1 and 3], typically a library of clones of chromosomal DNA of Streptococcus pneumoniae 0100993 in E.coli or some other suitable host is probed with a radiolabeled oligonucleotide, preferably a 17-mer or longer, derived from a partial sequence. Clones carrying DNA identical to that of the probe can then be distinguished using stringent conditions. By sequencing the individual clones thus identified with sequencing primers designed from the original sequence it is then possible to extend the sequence in both directions to determine the full gene sequence. Conveniently, such sequencing is performed using denatured double stranded DNA prepared from a plasmid clone. Suitable techniques are described by Maniatis, T., Fritsch, E. F. and Sambrook et al., MOLECULAR CLONING, A LABORATORY MANUAL, 2nd Ed.; Cold Spring Harbor Laboratory Press, Cold Spring Harbor, N.Y. (1989). (see in particular Screening By Hybridization 1.90 and Sequencing Denatured Double-Stranded DNA Templates 13.70). Illustrative of the invention, the polynucleotide set out in Table 1 [SEQ ID NO:1] was discovered in a DNA library derived from Streptococcus pneumoniae 0100993.\nThe DNA sequence set out in Table 1 [SEQ ID NOS:1] contains an open reading frame encoding a protein having about the number of amino acid residues set forth in Table 1 [SEQ ID NO:2] with a deduced molecular weight that can be calculated using amino acid residue molecular weight values well known in the art. The polynucleotide of SEQ ID NO: 1, between nucleotide number 1 through number 1023 encodes the polypeptide of SEQ ID NO:2. The stop codon begins at nucleotide number 1024 of SEQ ID NO:1.\ntrpS of the invention is structurally related to other proteins of the tryptophanyl tRNA synthetase family, as shown by the results of sequencing the DNA encoding trpS of the deposited strain. The protein exhibits greatest homology to Clostridium longisporum tryptophanyl tRNA synthetase protein among known proteins. trpS of Table 1 [SEQ ID NO:2] has about 66% identity over its entire length and about 81% similarity over its entire length with the amino acid sequence of Clostridium longisporum tryptophanyl tRNA synthetase polypeptide.\nThe invention provides a polynucleotide sequence identical over its entire length to the coding sequence in Table 1 [SEQ ID NO:1]. Also provided by the invention is the coding sequence for the mature polypeptide or a fragment thereof, by itself as well as the coding sequence for the mature polypeptide or a fragment in reading frame with other coding sequence, such as those encoding a leader or secretory sequence, a pre-, or pro- or prepro-protein sequence. The polynucleotide may also contain non-coding sequences, including for example, but not limited to non-coding 5xe2x80x2 and 3xe2x80x2 sequences, such as the transcribed, non-translated sequences, termination signals, ribosome binding sites, sequences that stabilize mRNA, introns, polyadenylation signals, and additional coding sequence which encode additional amino acids. For example, a marker sequence that facilitates purification of the fused polypeptide can be encoded. In certain embodiments of the invention, the marker sequence is a hexa-histidine peptide, as provided in the pQE vector (Qiagen, Inc.) and described in Gentz et al., Proc. Natl. Acad. Sci., USA 86: 821-824 (1989), or an HA tag (Wilson et al., Cell 37: 767 (1984). Polynucleotides of the invention also include, but are not limited to, polynucleotides comprising a structural gene and its naturally associated sequences that control gene expression.\nA preferred embodiment of the invention is a polynucleotide of comprising nucleotide 1 to 1023 or 1024 set forth in SEQ ID NO:1 of Table 1 which encode the trpS polypeptide.\nThe invention also includes polynucleotides of the formula set forth in Table 1 (C)[SEQ ID NO:1] wherein, at the 5xe2x80x2 end of the molecule, X is hydrogen, and at the 3xe2x80x2 end of the molecule, Y is hydrogen or a metal, R1 and R2 is any nucleic acid residue, and n is an integer between 1 and 1000. Any stretch of nucleic acid residues denoted by either R group, where R is greater than 1, may be either a heteropolymer or a homopolymer, preferably a heteropolymer.\nThe term xe2x80x9cpolynucleotide encoding a polypeptidexe2x80x9d as used herein encompasses polynucleotides that include a sequence encoding a polypeptide of the invention, particularly a bacterial polypeptide and more particularly a polypeptide of the Streptococcus pneumoniae trpS having the amino acid sequence set out in Table 1 [SEQ ID NO:2]. The term also encompasses polynucleotides that include a single continuous region or discontinuous regions encoding the polypeptide (for example, interrupted by integrated phage or an insertion sequence or editing) together with additional regions, that also may contain coding and/or non-coding sequences.\nThe invention further relates to variants of the polynucleotides described herein that encode for variants of the polypeptide having the deduced amino acid sequence of Table 1 [SEQ ID NO:2]. Variants that are fragments of the polynucleotides of the invention may be used to synthesize full-length polynucleotides of the invention.\nFurther particularly preferred embodiments are polynucleotides encoding trpS variants, that have the amino acid sequence of trpS polypeptide of Table 1 [SEQ ID NO:2] in which several, a few, 5 to 10, 1 to 5, 1 to 3, 2, 1 or no amino acid residues are substituted, deleted or added, in any combination. Especially preferred among these are silent substitutions, additions and deletions, that do not alter the properties and activities of trpS.\nFurther preferred embodiments of the invention are polynucleotides that are at least 70% identical over their entire length to a polynucleotide encoding trpS polypeptide having an amino acid sequence set out in Table 1 [SEQ ID NOS:2 and 4], and polynucleotides that are complementary to such polynucleotides. Alternatively, most highly preferred are polynucleotides that comprise a region that is at least 80% identical over its entire length to a polynucleotide encoding trpS polypeptide of the deposited strain and polynucleotides complementary thereto. In this regard, polynucleotides at least 90% identical over their entire length to the same are particularly preferred, and among these particularly preferred polynucleotides, those with at least 95% are especially preferred. Furthermore, those with at least 97% are highly preferred among those with at least 95%, and among these those with at least 98% and at least 99% are particularly highly preferred, with at least 99% being the more preferred.\nPreferred embodiment are polynucleotides that encode polypeptides that retain substantially the same biological function or activity as the mature polypeptide encoded by the DNA of Table 1 [SEQ ID NO:1].\nThe invention further relates to polynucleotides that hybridize to the herein above-described sequences. In this regard, the invention especially relates to polynucleotides that hybridize under stringent conditions to the herein above-described polynucleotides. As herein used, the terms xe2x80x9cstringent conditionsxe2x80x9d and xe2x80x9cstringent hybridization conditionsxe2x80x9d mean hybridization will occur only if there is at least 95% and preferably at least 97% identity between the sequences. An example of stringent hybridization conditions is overnight incubation at 42xc2x0 C. in a solution comprising: 50% formamide, 5xc3x97SSC (150 mM NaCl, 15 mM trisodium citrate), 50 mM sodium phosphate (pH7.6), 5xc3x97Denhardt\"\"s solution, 10% dextran sulfate, and 20 micrograms/ml denatured, sheared salmon sperm DNA, followed by washing the hybridization support in 0.1xc3x97SSC at about 65xc2x0 C. Hybridization and wash conditions are well known and exemplified in Sambrook, et al., Molecular Cloning: A Laboratory Manual, Second Edition, Cold Spring Harbor, N.Y., (1989), particularly Chapter 11 therein.\nThe invention also provides a polynucleotide consisting essentially of a polynucleotide sequence obtainable by screening an appropriate library containing the complete gene for a polynucleotide sequence set forth in SEQ ID NO:1 or SEQ ID NO:3 under stringent hybridization conditions with a probe having the sequence of said polynucleotide sequence set forth in SEQ ID NO:1 or a fragment thereof; and isolating said DNA sequence. Fragments useful for obtaining such a polynucleotide include, for example, probes and primers described elsewhere herein.\nAs discussed additionally herein regarding polynucleotide assays of the invention, for instance, polynucleotides of the invention as discussed above, may be used as a hybridization probe for RNA, cDNA and genomic DNA to isolate full-length cDNAs and genomic clones encoding trpS and to isolate cDNA and genomic clones of other genes that have a high sequence similarity to the trpS gene. Such probes generally will comprise at least 15 bases. Preferably, such probes will have at least 30 bases and may have at least 50 bases. Particularly preferred probes will have at least 30 bases and will have 50 bases or less.\nFor example, the coding region of the trpS gene may be isolated by screening using the DNA sequence provided in SEQ ID NO: 1 to synthesize an oligonucleotide probe. A labeled oligonucleotide having a sequence complementary to that of a gene of the invention is then used to screen a library of cDNA, genomic DNA or mRNA to determine which members of the library the probe hybridizes to.\nThe polynucleotides and polypeptides of the invention may be employed, for example, as research reagents and materials for discovery of treatments of and diagnostics for disease, particularly human disease, as further discussed herein relating to polynucleotide assays.\nPolynucleotides of the invention that are oligonucleotides derived from the sequences of SEQ ID NOS:1 and/or 2 may be used in the processes herein as described, but preferably for PCR, to determine whether or not the polynucleotides identified herein in whole or in part are transcribed in bacteria in infected tissue. It is recognized that such sequences will also have utility in diagnosis of the stage of infection and type of infection the pathogen has attained.\nThe invention also provides polynucleotides that may encode a polypeptide that is the mature protein plus additional amino or carboxyl-terminal amino acids, or amino acids interior to the mature polypeptide (when the mature form has more than one polypeptide chain, for instance). Such sequences may play a role in processing of a protein from precursor to a mature form, may allow protein transport, may lengthen or shorten protein half-life or may facilitate manipulation of a protein for assay or production, among other things. As generally is the case in vivo, the additional amino acids may be processed away from the mature protein by cellular enzymes.\nA precursor protein, having the mature form of the polypeptide fused to one or more prosequences may be an inactive form of the polypeptide. When prosequences are removed such inactive precursors generally are activated. Some or all of the prosequences may be removed before activation. Generally, such precursors are called proproteins.\nIn sum, a polynucleotide of the invention may encode a mature protein, a mature protein plus a leader sequence (which may be referred to as a preprotein), a precursor of a mature protein having one or more prosequences that are not the leader sequences of a preprotein, or a preproprotein, which is a precursor to a proprotein, having a leader sequence and one or more prosequences, which generally are removed during processing steps that produce active and mature forms of the polypeptide.\nVectors, Host Cells, Expression\nThe invention also relates to vectors that comprise a polynucleotide or polynucleotides of the invention, host cells that are genetically engineered with vectors of the invention and the production of polypeptides of the invention by recombinant techniques. Cell-free translation systems can also be employed to produce such proteins using RNAs derived from the DNA constructs of the invention.\nFor recombinant production, host cells can be genetically engineered to incorporate expression systems or portions thereof or polynucleotides of the invention. Introduction of a polynucleotide into the host cell can be effected by methods described in many standard laboratory manuals, such as Davis et al., BASIC METHODS IN MOLECULAR BIOLOGY, (1986) and Sambrook et al., MOLECULAR CLONING: A LABORATORY MANUAL, 2nd Ed., Cold Spring Harbor Laboratory Press, Cold Spring Harbor, N.Y. (1989), such as, calcium phosphate transfection, DEAE-dextran mediated transfection, transvection, microinjection, cationic lipid-mediated transfection, electroporation, transduction, scrape loading, ballistic introduction and infection.\nRepresentative examples of appropriate hosts include bacterial cells, such as streptococci, staphylococci, enterococci E. coli , streptomyces and Bacillus subtilis cells; fungal cells, such as yeast cells and Aspergillus cells; insect cells such as Drosophila S2 and Spodoptera Sf9 cells; animal cells such as CHO, COS, HeLa, C127, 3T3, BHK, 293 and Bowes melanoma cells; and plant cells.\nA great variety of expression systems can be used to produce the polypeptides of the invention. Such vectors include, among others, chromosomal, episomal and virus-derived vectors, e.g., vectors derived from bacterial plasmids, from bacteriophage, from transposons, from yeast episomes, from insertion elements, from yeast chromosomal elements, from viruses such as baculoviruses, papova viruses, such as SV40, vaccinia viruses, adenoviruses, fowl pox viruses, pseudorabies viruses and retroviruses, and vectors derived from combinations thereof, such as those derived from plasmid and bacteriophage genetic elements, such as cosmids and phagemids. The expression system constructs may contain control regions that regulate as well as engender expression. Generally, any system or vector suitable to maintain, propagate or express polynucleotides and/or to express a polypeptide in a host may be used for expression in this regard. The appropriate DNA sequence may be inserted into the expression system by any of a variety of well-known and routine techniques, such as, for example, those set forth in Sambrook et al., MOLECULAR CLONING, A LABORATORY MANUAL, (supra).\nFor secretion of the translated protein into the lumen of the endoplasmic reticulum, into the periplasmic space or into the extracellular environment, appropriate secretion signals may be incorporated into the expressed polypeptide. These signals may be endogenous to the polypeptide or they may be heterologous signals.\nPolypeptides of the invention can be recovered and purified from recombinant cell cultures by well-known methods including ammonium sulfate or ethanol precipitation, acid extraction, anion or cation exchange chromatography, phosphocellulose chromatography, hydrophobic interaction chromatography, affinity chromatography, hydroxylapatite chromatography, and lectin chromatography. Most preferably, high performance liquid chromatography is employed for purification. Well known techniques for refolding protein may be employed to regenerate active conformation when the polypeptide is denatured during isolation and or purification.\nDiagnostic Assays\nThis invention is also related to the use of the trpS polynucleotides of the invention for use as diagnostic reagents. Detection of trpS in a eukaryote, particularly a mammal, and especially a human, will provide a diagnostic method for diagnosis of a disease. Eukaryotes (herein also xe2x80x9cindividual(s)xe2x80x9d), particularly mammals, and especially humans, particularly those infected or suspected to be infected with an organism comprising the trpS gene may be detected at the nucleic acid level by a variety of techniques.\nNucleic acids for diagnosis may be obtained from an infected individual\"\"s cells and tissues, such as bone, blood, muscle, cartilage, and skin. Genomic DNA may be used directly for detection or may be amplified enzymatically by using PCR or other amplification technique prior to analysis. RNA or cDNA may also be used in the same ways. Using amplification, characterization of the species and strain of prokaryote present in an individual, may be made by an analysis of the genotype of the prokaryote gene. Deletions and insertions can be detected by a change in size of the amplified product in comparison to the genotype of a reference sequence. Point mutations can be identified by hybridizing amplified DNA to labeled trpS polynucleotide sequences. Perfectly matched sequences can be distinguished from mismatched duplexes by RNase digestion or by differences in melting temperatures. DNA sequence differences may also be detected by alterations in the electrophoretic mobility of the DNA fragments in gels, with or without denaturing agents, or by direct DNA sequencing. See, e.g., Myers et al., Science, 230: 1242 (1985). Sequence changes at specific locations also may be revealed by nuclease protection assays, such as RNase and S1 protection or a chemical cleavage method. See, e.g., Cotton et al., Proc. Natl. Acad Sci., USA, 85: 4397-4401 (1985).\nCells carrying mutations or polymorphisms in the gene of the invention may also be detected at the DNA level by a variety of techniques, to allow for serotyping, for example. For example, RT-PCR can be used to detect mutations. It is particularly preferred to used RT-PCR in conjunction with automated detection systems, such as, for example, GeneScan. RNA or cDNA may also be used for the same purpose, PCR or RT-PCR. As an example, PCR primers complementary to a nucleic acid encoding trpS can be used to identify and analyze mutations. Examples of representative primers are shown below in Table 2.\nThe invention further provides these primers with 1, 2, 3 or 4 nucleotides removed from the 5xe2x80x2 and/or the 3xe2x80x2 end. These primers may be used for, among other things, amplifying trpS DNA isolated from a sample derived from an individual. The primers may be used to amplify the gene isolated from an infected individual such that the gene may then be subject to various techniques for elucidation of the DNA sequence. In this way, mutations in the DNA sequence may be detected and used to diagnose infection and to serotype and/or classify the infectious agent.\nThe invention further provides a process for diagnosing, disease, preferably bacterial infections, more preferably infections by Streptococcus pneumoniae, and most preferably otitis media, conjunctivitis, pneumonia, bacteremia, meningitis, sinusitis, pleural empyema and endocarditis, and most particularly meningitis, such as for example infection of cerebrospinal fluid, comprising determining from a sample derived from an individual a increased level of expression of polynucleotide having the sequence of Table 1 [SEQ ID NO: 1]. Increased or decreased expression of trpS polynucleotide can be measured using any on of the methods well known in the art for the quantation of polynucleotides, such as, for example, amplification, PCR, RT-PCR, RNase protection, Northern blotting and other hybridization methods.\nIn addition, a diagnostic assay in accordance with the invention for detecting over-expression of trpS protein compared to normal control tissue samples may be used to detect the presence of an infection, for example. Assay techniques that can be used to determine levels of a trpS protein, in a sample derived from a host are well-known to those of skill in the art. Such assay methods include radioimmunoassays, competitive-binding assays, Western Blot analysis and ELISA assays.\nAntibodies\nThe polypeptides of the invention or variants thereof, or cells expressing them can be used as an immunogen to produce antibodies immunospecific for such polypeptides. xe2x80x9cAntibodiesxe2x80x9d as used herein includes monoclonal and polyclonal antibodies, chimeric, single chain, simianized antibodies and humanized antibodies, as well as Fab fragments, including the products of an Fab immunoglobulin expression library.\nAntibodies generated against the polypeptides of the invention can be obtained by administering the polypeptides or epitope-bearing fragments, analogues or cells to an animal, preferably a nonhuman, using routine protocols. For preparation of monoclonal antibodies, any technique known in the art that provides antibodies produced by continuous cell line cultures can be used. Examples include various techniques, such as those in Kohler, G. and Milstein, C., Nature 256: 495-497 (1975); Kozbor et al., Immunology Today 4: 72 (1983); Cole et al., pg. 77-96 in MONOCLONAL ANTIBODIES AND CANCER THERAPY, Alan R. Liss, Inc. (1985).\nTechniques for the production of single chain antibodies (U.S. Pat. No. 4,946,778) can be adapted to produce single chain antibodies to polypeptides of this invention. Also, transgenic mice, or other organisms such as other mammals, may be used to express humanized antibodies.\nAlternatively phage display technology may be utilized to select antibody genes with binding activities towards the polypeptide either from repertoires of PCR amplified v-genes of lymphocytes from humans screened for possessing anti-trpS or from naive libraries (McCafferty, J. et al., (1990), Nature 348, 552-554; Marks, J. et al., (1992) Biotechnology 10, 779-783). The affinity of these antibodies can also be improved by chain shuffling (Clackson, T. et al., (1991) Nature 352, 624-628).\nIf two antigen binding domains are present each domain may be directed against a different epitopexe2x80x94termed xe2x80x98bispecificxe2x80x99 antibodies.\nThe above-described antibodies may be employed to isolate or to identify clones expressing the polypeptides to purify the polypeptides by affinity chromatography.\nThus, among others, antibodies against trpS-polypeptide may be employed to treat infections, particularly bacterial infections and especially otitis media, conjunctivitis, pneumonia, bacteremia, meningitis, sinusitis, pleural empyema and endocarditis, and most particularly meningitis, such as for example infection of cerebrospinal fluid.\nPolypeptide variants include antigenically, epitopically or immunologically equivalent variants that form a particular aspect of this invention. The term xe2x80x9cantigenically equivalent derivativexe2x80x9d as used herein encompasses a polypeptide or its equivalent which will be specifically recognized by certain antibodies which, when raised to the protein or polypeptide according to the invention, interfere with the immediate physical interaction between pathogen and mammalian host. The term xe2x80x9cimmunologically equivalent derivativexe2x80x9d as used herein encompasses a peptide or its equivalent which when used in a suitable formulation to raise antibodies in a vertebrate, the antibodies act to interfere with the immediate physical interaction between pathogen and mammalian host.\nThe polypeptide, such as an antigenically or immunologically equivalent derivative or a fusion protein thereof is used as an antigen to immunize a mouse or other animal such as a rat or chicken. The fusion protein may provide stability to the polypeptide. The antigen may be associated, for example by conjugation, with an immunogenic carrier protein for example bovine serum albumin (BSA) or keyhole limpet haemocyanin (KLH). Alternatively a multiple antigenic peptide comprising multiple copies of the protein or polypeptide, or an antigenically or immunologically equivalent polypeptide thereof may be sufficiently antigenic to improve immunogenicity so as to obviate the use of a carrier.\nPreferably, the antibody or variant thereof is modified to make it less immunogenic in the individual. For example, if the individual is human the antibody may most preferably be xe2x80x9chumanizedxe2x80x9d; where the complimentarity determining region(s) of the hybridoma-derived antibody has been transplanted into a human monoclonal antibody, for example as described in Jones, P. et al. (1986), Nature 321, 522-525 or Tempest et al., (1991) Biotechnology 9, 266-273.\nThe use of a polynucleotide of the invention in genetic immunization will preferably employ a suitable delivery method such as direct injection of plasmid DNA into muscles (Wolff et al., Hum Mol Genet 1992, 1:363, Manthorpe et al., Hum. Gene Ther. 1963:4, 419), delivery of DNA complexed with specific protein carriers (Wu et al., J Biol Chem. 1989: 264,16985), coprecipitation of DNA with calcium phosphate (Benvenisty and Reshef, PNAS USA, 1986:83,9551), encapsulation of DNA in various forms of liposomes (Kaneda et al., Science 1989:243,375), particle bombardment (Tang et al., Nature 1992, 356:152, Eisenbraun et al., DNA Cell Biol 1993, 12:791) and in vivo infection using cloned retroviral vectors (Seeger et al., PNAS USA 1984:81,5849).\nAntagonists and Agonistsxe2x80x94Assays and Molecules\nPolypeptides of the invention may also be used to assess the binding of small molecule substrates and ligands in, for example, cells, cell-free preparations, chemical libraries, and natural product mixtures. These substrates and ligands may be natural substrates and ligands or may be structural or functional mimetics. See, e.g., Coligan et al., Current Protocols in Immunology 1(2): Chapter 5 (1991).\nThe invention also provides a method of screening compounds to identify those which enhance (agonist) or block (antagonist) the action of trpS polypeptides or polynucleotides, particularly those compounds that are bacteriostatic and/or bacteriocidal. The method of screening may involve high-throughput techniques. For example, to screen for agonists or antagonists, a synthetic reaction mix, a cellular compartment, such as a membrane, cell envelope or cell wall, or a preparation of any thereof, comprising trpS polypeptide and a labeled substrate or ligand of such polypeptide is incubated in the absence or the presence of a candidate molecule that may be a trpS agonist or antagonist. The ability of the candidate molecule to agonize or antagonize the trpS polypeptide is reflected in decreased binding of the labeled ligand or decreased production of product from such substrate. Molecules that bind gratuitously, i.e., without inducing the effects of trpS polypeptide are most likely to be good antagonists. Molecules that bind well and increase the rate of product production from substrate are agonists. Detection of the rate or level of production of product from substrate may be enhanced by using a reporter system. Reporter systems that may be useful in this regard include but are not limited to colorimetric labeled substrate converted into product, a reporter gene that is responsive to changes in trpS polynucleotide or polypeptide activity, and binding assays known in the art.\nAnother example of an assay for trpS antagonists is a competitive assay that combines trpS and a potential antagonist with trpS-binding molecules, recombinant trpS binding molecules, natural substrates or ligands, or substrate or ligand mimetics, under appropriate conditions for a competitive inhibition assay. trpS can be labeled, such as by radioactivity or a colorimetric compound, such that the number of trpS molecules bound to a binding molecule or converted to product can be determined accurately to assess the effectiveness of the potential antagonist.\nPotential antagonists include small organic molecules, peptides, polypeptides and antibodies that bind to a polynucleotide or polypeptide of the invention and thereby inhibit or extinguish its activity. Potential antagonists also may be small organic molecules, a peptide, a polypeptide such as a closely related protein or antibody that binds the same sites on a binding molecule, such as a binding molecule, without inducing trpS-induced activities, thereby preventing the action of trpS by excluding trpS from binding.\nPotential antagonists include a small molecule that binds to and occupies the binding site of the polypeptide thereby preventing binding to cellular binding molecules, such that normal biological activity is prevented. Examples of small molecules include but are not limited to small organic molecules, peptides or peptide-like molecules. Other potential antagonists include antisense molecules (see Okano, J. Neurochem. 56: 560 (1991); OLIGODEOXYNUCLEOTIDES AS ANTISENSE INHIBITORS OF GENE EXPRESSION, CRC Press, Boca Raton, Fla. (1988), for a description of these molecules). Preferred potential antagonists include compounds related to and variants of trpS.\nEach of the DNA sequences provided herein may be used in the discovery and development of antibacterial compounds. The encoded protein, upon expression, can be used as a target for the screening of antibacterial drugs. Additionally, the DNA sequences encoding the amino terminal regions of the encoded protein or Shine-Delgarno or other translation facilitating sequences of the respective mRNA can be used to construct antisense sequences to control the expression of the coding sequence of interest.\nThe invention also provides the use of the polypeptide, polynucleotide or inhibitor of the invention to interfere with the initial physical interaction between a pathogen and mammalian host responsible for sequelae of infection. In particular the molecules of the invention may be used: in the prevention of adhesion of bacteria, in particular gram positive bacteria, to mammalian extracellular matrix proteins on in-dwelling devices or to extracellular matrix proteins in wounds; to block trpS protein-mediated mammalian cell invasion by, for example, initiating phosphorylation of mammalian tyrosine kinases (Rosenshine et al., Infect. Immun. 60:2211 (1992); to block bacterial adhesion between mammalian extracellular matrix proteins and bacterial trpS proteins that mediate tissue damage and; to block the normal progression of pathogenesis in infections initiated other than by the implantation of in-dwelling devices or by other surgical techniques.\nThe antagonists and agonists of the invention may be employed, for instance, to inhibit and treat otitis media, conjunctivitis, pneumonia, bacteremia, meningitis, sinusitis, pleural empyema and endocarditis, and most particularly meningitis, such as for example infection of cerebrospinal fluid.\nHelicobacter pylori (herein H. pylori) bacteria infect the stomachs of over one-third of the world\"\"s population causing stomach cancer, ulcers, and gastritis (International Agency for Research on Cancer (1994) Schistosomes, Liver Flukes and Helicobacter Pylori (International Agency for Research on Cancer, Lyon, France; http://www.uicc.ch/ecp/ecp2904.htm). Moreover, the international Agency for Research on Cancer recently recognized a cause-and-effect relationship between H. pylori and gastric adenocarcinoma, classifying the bacterium as a Group I (definite) carcinogen. Preferred antimicrobial compounds of the invention (agonists and antagonists of trpS) found using screens provided by the invention, particularly broad-spectrum antibiotics, should be useful in the treatment of H. pylori infection. Such treatment should decrease the advent of H. pylori-induced cancers, such as gastrointestinal carcinoma. Such treatment should also cure gastric ulcers and gastritis.\nVaccines\nAnother aspect of the invention relates to a method for inducing an immunological response in an individual, particularly a mammal which comprises inoculating the individual with trpS, or a fragment or variant thereof, adequate to produce antibody and/ or T cell immune response to protect said individual from infection, particularly bacterial infection and most particularly Streptococcus pneumoniae infection. Also provided are methods whereby such immunological response slows bacterial replication. Yet another aspect of the invention relates to a method of inducing immunological response in an individual which comprises delivering to such individual a nucleic acid vector to direct expression of trpS, or a fragment or a variant thereof, for expressing trpS, or a fragment or a variant thereof in vivo in order to induce an immunological response, such as, to produce antibody and/ or T cell immune response, including, for example, cytokine-producing T cells or cytotoxic T cells, to protect said individual from disease, whether that disease is already established within the individual or not. One way of administering the gene is by accelerating it into the desired cells as a coating on particles or otherwise. Such nucleic acid vector may comprise DNA, RNA, a modified nucleic acid, or a DNA/RNA hybrid.\nA further aspect of the invention relates to an immunological composition which, when introduced into an individual capable or having induced within it an immunological response, induces an immunological response in such individual to a trpS or protein coded therefrom, wherein the composition comprises a recombinant trpS or protein coded therefrom comprising DNA which codes for and expresses an antigen of said trpS or protein coded therefrom. The immunological response may be used therapeutically or prophylactically and may take the form of antibody immunity or cellular immunity such as that arising from CTL or CD4+ T cells.\nA trpS polypeptide or a fragment thereof may be fused with co-protein which may not by itself produce antibodies, but is capable of stabilizing the first protein and producing a fused protein which will have immunogenic and protective properties. Thus fused recombinant protein, preferably further comprises an antigenic co-protein, such as lipoprotein D from Hemophilus influenzae, Glutathione-S-transferase (GST) or beta-galactosidase, relatively large co-proteins which solubilize the protein and facilitate production and purification thereof. Moreover, the co-protein may act as an adjuvant in the sense of providing a generalized stimulation of the immune system. The co-protein may be attached to either the amino or carboxy terminus of the first protein.\nProvided by this invention are compositions, particularly vaccine compositions, and methods comprising the polypeptides or polynucleotides of the invention and immunostimulatory DNA sequences, such as those described in Sato, Y. et al. Science 273: 352 (1996).\nAlso, provided by this invention are methods using the described polynucleotide or particular fragments thereof which have been shown to encode non-variable regions of bacterial cell surface proteins in DNA constructs used in such genetic immunization experiments in animal models of infection with Streptococcus pneumoniae will be particularly useful for identifying protein epitopes able to provoke a prophylactic or therapeutic immune response. It is believed that this approach will allow for the subsequent preparation of monoclonal antibodies of particular value from the requisite organ of the animal successfully resisting or clearing infection for the development of prophylactic agents or therapeutic treatments of bacterial infection, particularly Streptococcus pneumoniae infection, in mammals, particularly humans.\nThe polypeptide may be used as an antigen for vaccination of a host to produce specific antibodies which protect against invasion of bacteria, for example by blocking adherence of bacteria to damaged tissue. Examples of tissue damage include wounds in skin or connective tissue caused, e.g., by mechanical, chemical or thermal damage or by implantation of indwelling devices, or wounds in the mucous membranes, such as the mouth, mammary glands, urethra or vagina.\nThe invention also includes a vaccine formulation which comprises an immunogenic recombinant protein of the invention together with a suitable carrier. Since the protein may be broken down in the stomach, it is preferably administered parenterally, including, for example, administration that is subcutaneous, intramuscular, intravenous, or intradermal. Formulations suitable for parenteral administration include aqueous and non-aqueous sterile injection solutions which may contain anti-oxidants, buffers, bacteriostats and solutes which render the formulation insotonic with the bodily fluid, preferably the blood, of the individual; and aqueous and non-aqueous sterile suspensions which may include suspending agents or thickening agents. The formulations may be presented in unit-dose or multi-dose containers, for example, sealed ampules and vials and may be stored in a freeze-dried condition requiring only the addition of the sterile liquid carrier immediately prior to use. The vaccine formulation may also include adjuvant systems for enhancing the immunogenicity of the formulation, such as oil-in water systems and other systems known in the art. The dosage will depend on the specific activity of the vaccine and can be readily determined by routine experimentation.\nWhile the invention has been described with reference to certain trpS protein, it is to be understood that this covers fragments of the naturally occurring protein and similar proteins with additions, deletions or substitutions which do not substantially affect the immunogenic properties of the recombinant protein.\nCompositions, Kits and Administration\nThe invention also relates to compositions comprising the polynucleotide or the polypeptides discussed above or their agonists or antagonists. The polypeptides of the invention may be employed in combination with a non-sterile or sterile carrier or carriers for use with cells, tissues or organisms, such as a pharmaceutical carrier suitable for administration to a subject. Such compositions comprise, for instance, a media additive or a therapeutically effective amount of a polypeptide of the invention and a pharmaceutically acceptable carrier or excipient. Such carriers may include, but are not limited to, saline, buffered saline, dextrose, water, glycerol, ethanol and combinations thereof. The formulation should suit the mode of administration. The invention further relates to diagnostic and pharmaceutical packs and kits comprising one or more containers filled with one or more of the ingredients of the aforementioned compositions of the invention.\nPolypeptides and other compounds of the invention may be employed alone or in conjunction with other compounds, such as therapeutic compounds.\nThe pharmaceutical compositions may be administered in any effective, convenient manner including, for instance, administration by topical, oral, anal, vaginal, intravenous, intraperitoneal, intramuscular, subcutaneous, intranasal or intradermal routes among others.\nIn therapy or as a prophylactic, the active agent may be administered to an individual as an injectable composition, for example as a sterile aqueous dispersion, preferably isotonic.\nAlternatively the composition may be formulated for topical application for example in the form of ointments, creams, lotions, eye ointments, eye drops, ear drops, mouthwash, impregnated dressings and sutures and aerosols, and may contain appropriate conventional additives, including, for example, preservatives, solvents to assist drug penetration, and emollients in ointments and creams. Such topical formulations may also contain compatible conventional carriers, for example cream or ointment bases, and ethanol or oleyl alcohol for lotions. Such carriers may constitute from about 1% to about 98% by weight of the formulation; more usually they will constitute up to about 80% by weight of the formulation.\nFor administration to mammals, and particularly humans, it is expected that the daily dosage level of the active agent will be from 0.01 mg/kg to 10 mg/kg, typically around 1 mg/kg. The physician in any event will determine the actual dosage which will be most suitable for an individual and will vary with the age, weight and response of the particular individual. The above dosages are exemplary of the average case. There can, of course, be individual instances where higher or lower dosage ranges are merited, and such are within the scope of this invention.\nIn-dwelling devices include surgical implants, prosthetic devices and catheters, i.e., devices that are introduced to the body of an individual and remain in position for an extended time. Such devices include, for example, artificial joints, heart valves, pacemakers, vascular grafts, vascular catheters, cerebrospinal fluid shunts, urinary catheters, continuous ambulatory peritoneal dialysis (CAPD) catheters.\nThe composition of the invention may be administered by injection to achieve a systemic effect against relevant bacteria shortly before insertion of an in-dwelling device. Treatment may be continued after surgery during the in-body time of the device. In addition, the composition could also be used to broaden perioperative cover for any surgical technique to prevent bacterial wound infections, especially Streptococcus pneumoniae wound infections.\nMany orthopaedic surgeons consider that humans with prosthetic joints should be considered for antibiotic prophylaxis before dental treatment that could produce a bacteremia. Late deep infection is a serious complication sometimes leading to loss of the prosthetic joint and is accompanied by significant morbidity and mortality. It may therefore be possible to extend the use of the active agent as a replacement for prophylactic antibiotics in this situation.\nIn addition to the therapy described above, the compositions of this invention may be used generally as a wound treatment agent to prevent adhesion of bacteria to matrix proteins exposed in wound tissue and for prophylactic use in dental treatment as an alternative to, or in conjunction with, antibiotic prophylaxis.\nAlternatively, the composition of the invention may be used to bathe an indwelling device immediately before insertion. The active agent will preferably be present at a concentration of 1 xcexcg/ml to 10 mg/ml for bathing of wounds or indwelling devices.\nA vaccine composition is conveniently in injectable form. Conventional adjuvants may be employed to enhance the immune response. A suitable unit dose for vaccination is 0.5-5 microgram/kg of antigen, and such dose is preferably administered 1-3 times and with an interval of 1-3 weeks. With the indicated dose range, no adverse toxicological effects will be observed with the compounds of the invention which would preclude their administration to suitable individuals.\nEach reference disclosed herein is incorporated by reference herein in its entirety. Any patent application to which this application claims priority is also incorporated by reference herein in its entirety."}
{"text": "In a series of previous issued and pending U.S. Patent documents inventors named in connection with the present document and various groupings of their colleagues have disclosed a family of electronic warfare radio receivers that are usable in obtaining military and otherwise useful information from an incoming microwave frequency radio signal and have also disclosed significant components of such receivers. Early in this sequence of patent documents the name “monobit receiver” was used in referring to both one embodiment of such a receiver and subsequently as a generic or family name for such receivers. Each of these receivers as preferably embodied is of a wide band nature and employs the Fourier transformation as a tool for identifying characteristics of the incoming radio frequency signal. Additionally each of these receivers has the attribute of employing a simplified multiplication arrangement in making use of the Fourier transformation results in order that a simple and low cost realization of the receiver can be achieved.\nIn each of these simplified multiplications the Kernel function portion of the Fourier transformation is represented by a unit value, a magnitude of one or near one, in order to avoid the mechanization complexity of full fledged multiplication algorithm during the receiver step when complex numbers representing an input signal are multiplied by a second set of complex numbers representing the Fourier transformation Kernel function. Generally these simplified multiplications have in the past been accomplished by insuring the second set of complex numbers, i.e., the Fourier transformation results, involve real and imaginary values of unity or near unity. As implied by the name “monobit receiver” the earliest of these simplified Kernel function receivers employs a Kernel function approximation representable by a single binary bit of information while the later versions of the receiver and indeed the present invention involve Kernel function arrangements requiring a greater number of bits, i.e., three or four bits. In view of these greater numbers of bits it is of course necessary to accomplish an actual multiplication operation involving the Fourier transformation results however for feasibility purposes some form of simplified and readily accomplished multiplication is needed.\nIn order to better appreciate the present invention it may be helpful to regress briefly and consider the evolution of simplified Kernel function realizations in our work. Such evolution commences with the realization of a Fourier transformation using only multiplication by unity or in essence no multiplication in the Fourier transformation computation algorithm. Kernel function realization in this manner is first disclosed in a U.S. patent of Tsui et al., a patent numbered U.S. Pat. No. 5,917,737, a patent wherein Kernel function values are located on a circle of unit radius at angular locations of 5/4, 3 5/4, 5 5/4 and 7 5/4 radians i.e., at locations displaced by forty five degrees from the real and imaginary axes of a coordinate axis plot. Actual coordinate axis locations of 1+j, 1−j, −1+j and −1+j are used for the Kernel functions disclosed in the U.S. Pat. No. 5,917,737 patent. For reference simplification purposes this and the several other patents issued to various combinations of persons from our laboratory are herein referred-to as “our patent” notwithstanding the differing group of inventors, with perhaps one common inventor, associated with each patent.\nOur U.S. Pat. No. 5,793,323 therefore relates to the U.S. Pat. No. 5,917,737 patent in that it discloses a single integrated circuit chip arrangement for a monobit receiver employing the approximated Kernel function of the U.S. Pat. No. 5,917,737 patent. This patent also discloses several implementation compromises possible in applying the simplified Kernel function to receiver apparatus and provides a receiver that can cover a one gigahertz spectrum with a frequency resolution of 9.77 megahertz.\nOur U.S. patent application Ser. No. 09/917,589 filed on Jul. 30, 2001 also relates to the U.S. Pat. No. 5,917,737 patent in that it discloses the use of a monobit receiver in combination with several bandpass filters in order to increase the number of simultaneously processable signals and enhance the dynamic range capability of the overall system.\nThe invention of our U.S. patent application Ser. No. 09/944,616, filed on Sep. 4, 2001, provides a straightforward approach to the enhancement of dynamic range in a monobit family receiver by increasing the number of Kernel function locations used in the Kernel function approximation from the four locations of the U.S. Pat. No. 5,917,737 patent and the U.S. Pat. No. 5,963,164 patent to eight locations. In the application of Ser. No. 09/944,616 Kernel function values located at the 5/4, 3 5/4, 5 5/4 and 7 5/4 radian locations are added to the Kernel function values at 0, 5/2, 5, and 3 5/2 radians with the added four values being slightly increased in magnitude from true unit circle values and in fact having a magnitude of (2)1/2 or 1.414. Some of these eight location Kernel function values depart slightly from unity magnitude however the overall results of the eight-location Kernel function appear favorable.\nAdditional of our U.S. Patent and Trademark Office documents involving Kernel function realizations include the application of Ser. No. 10/115,819, filed on Apr. 3, 2002, and now issued as Statutory Invention Registration H2109. This document is titled “PASSIVE MICROWAVE DIRECTION FINDING WITH MONOBIT FOURIER TRANSFORMATION RECEIVER AND MATRIX COUPLED ANTENNA”\nIn addition to these previous patent documents wherein use of a Kernel function realization of unity or near unity value is employed in order to achieve a simplified multiplication operation the U.S. patent application Ser. No. 10/354,267, filed on Jan. 31, 2003, discloses an electronic warfare receiver wherein a Kernel function realization of other than unity magnitude is used in combination with a simplified multiplication arrangement accomplished with a bit shifting multiplication algorithm. In the of Ser. No. 10/354,267, document the simplified bit shifting algorithm multiplication can be implemented with such minimal complexity and cost as to make the non-unity magnitude of the Kernel function of little consequence. The present invention may be considered as an extension of this simplification multiplication through shifting concept.\nThe present invention may also start with the initial realization that a significant difficulty with the monobit electronic warfare receiver results form the fact that the two signal instantaneous dynamic range achieved in such a receiver is undesirably low. The instantaneous dynamic range relates to the receiver's capability to detect two simultaneous signals of different amplitude and can be an important consideration in a receiver use situation especially when a hostile adversary is involved. Because of the non-linear property of the approximated Kernel function in the Tsui et al. U.S. Pat. No. 5,917,737 patent receiver, the two signal instantaneous dynamic range of this receiver is limited to about 2 dB. Using a different Kernel function in the later tribit receiver of the Ser. No. 09/944,616 application can improve this instantaneous dynamic range to about 10 dB. A receiver arrangement according to the present invention can further improve this dynamic range characteristic. Each of the patent documents, issued or pending, identified in this BACKGROUND OF THE INVENTION DISCUSSION is hereby incorporated by reference herein.\nThe present patent document continues in this line of Kernel function and Kernel function receiver patents and thus provides additional insight into the simplified Kernel function utilization art. In the present invention additional improvement of the dynamic range achieved in the monobit family of receivers is pursued while relaxing the previously perceived need to limit multiplications to unity magnitude values."}
{"text": "Backside Illumination (BSI) image sensor chips are replacing front-side illumination sensor chips for their higher efficiency in capturing photons. In the formation of the BSI image sensor chips, image sensors and logic circuits are formed on a silicon substrate of a wafer, followed by the formation of an interconnect structure on a front side of the silicon chip.\nThe image sensors in the BSI image sensor chips generate electrical signals in response to the stimulation of photons. The magnitudes of the electrical signals (such as the currents) depend on the intensity of the incident light received by the respective image sensors. The image sensors, however, suffer from non-optically generated signals, which include the leakage signals, thermally generated signals, dark currents, and the like. Accordingly, the electrical signals generated by the image sensors need to be calibrated, so that the undesirable signals are cancelled out from the output signals of the image sensors. To cancel the non-optically generated signals, black reference image sensors are formed, and are used to generate non-optically generated signals. The black reference image sensors, therefore, need to be blocked from receiving light signals."}
{"text": "It has long been known in the art that a contact lens must have surfaces that have a certain degree of hydrophilicity in order to be wet by tears thus providing unblurred vision.\nSoft, hydrophilic contact lenses, in addition to being wettable, provide comfort to the wearer but lack the ability to correct certain visual deficiencies such as astigmatism since they tend to conform to the shape of the corneal surface.\nOften hydrophilic monomers can be added to a mixture of comonomers in the formation of contact lenses so that upon polymerization optically clear contact lenses result which have a certain degree of hydrophilicity. As the hydrophilic monomer content increases where it is added directly to the lens composition, the physical characteristics of the lenses are affected by the increased hydration propensity of the polymeric composition.\nIn some cases, it has been known to treat a formed contact lens with a polymerizable hydrophilic monomer to form a surface coating of hydrophilic polymer grafted to an otherwise hydrophobic polymer surface. Although effective, this method of increasing the hydrophilic character of the lens surface can suffer from involved and difficult manufacturing procedures as well as lack of permanence.\nPresent rigid and soft contact lenses sometimes retain water on their surfaces through secondary chemical bonding and as a consequence only a very thin layer of water molecules is present between the eye and the contact lens.\nSoft lenses are inherently comfortable but often times, as with hard lenses, suffer from brief surface dryness between eye blinks. State of the art technology teaches that a water soluble neutral polymer may be applied to the surfaces of a hard contact lens to provide a \"cushion\" layer between the lens and the eye which is equated with increased wettability as well as wearer comfort and tolerance.\nDissipation of the \"cushion\" layer occurs rapidly in most prior art constructions, since there is little specific interaction between the mobile polymer in this layer and the lens surface. As a result the wearer begins to feel discomfort and must rewet the lens surfaces."}
{"text": "Knowledge and information derivable from employees within an organization are often undervalued key resources. The risk in improperly managing these resources is usually underestimated and its impact not felt until the people holding the requisite knowledge to ensure that certain processes run smoothly within the organization leave. Sometimes, an organization may also be unaware of lost opportunities to benefit from the knowledge of its people, or for that matter knowledge in general, because of a lack of an open support system that encourages the sharing of expertise and know-how to other members within the organization.\nThere is clearly a need for a system for acquiring, sharing and managing knowledge and information within an organization."}
{"text": "1. Field of the Invention\nThe present invention relates to an electronic mail cataloging and retrieving system which is used to catalog and retrieve many electronic mails exchanged by a wide area network, and particularly to an electronic mail retrieving system which saves time and trouble required for cataloging and retrieving and facilitates the addition and change of electronic mail classifying rules.\n2. Description of the Related Art\nA conventional electronic mail cataloging and retrieving system of this type which is used for cataloging and retrieving electronic mails catalogues and retrieves lots of electronic mails by utilizing a virtually set mail folder. The mail folder here is univocally given a meaning by a word or word string, and each mail folder has electronic mails conforming to a set meaning in it. Namely, electronic mails are classified by a plurality of mail folders and their meanings. And, when a mail folder is opened, electronic mails kept in it can be seen. On the other hand, when a mail folder is in a closed state, electronic mails kept in the mail folder cannot be seen.\nThe conventional electronic mail cataloging and retrieving system uses the above mail folders to allow free cataloging and retrieving of electronic mails sent or received by a wide area network according to a user's taste.\nFIG. 39 shows one example of the screen shown by a conventional electronic mail cataloging and retrieving system. The conventional electronic mail cataloging and retrieving system will be described with reference to FIG. 39.\nAs shown in the drawing, a screen 2900 has a mail folder display area 2901 which shows a list of mail folders defined within the electronic mail cataloging and retrieving system and a header information display area 2902 which shows a list of header information on electronic mails. Electronic mails of which the list of header information is shown on the header information screen 2902 are those stored in an open mail folder.\nThe mail folder display area 2901 shows a mail folder state, namely a folder symbol for showing whether or not the mail folder is open, and a folder name for showing the meaning given to the mail folder. In FIG. 39, the folder symbol corresponding to the folder name indicates that a mail folder having a folder name \"Inbox\" for example is open.\nThe header information display area 2902 shows a sending date and time, sender, title and text in part (behind the symbol &lt;&lt; in the drawing) extracted from an electronic mail stored in the open mail folder in the form modified to an at-a-glance guide.\nWhen another mail folder other than the open mail folder is to be opened by the electronic mail cataloging and retrieving system, the following operation and process (1) to (3) are performed.\n(1) One of the mail folders on which information is shown in the mail folder display area 2901 is selected.\n(2) The folder symbol of the mail folder selected by the operation (1) is changed to a \"folder symbol indicating its folder is open\".\n(3) The header information display area 2902 shows the header information on electronic mails stored in the mail folder newly selected by the operation (1) instead of the header information on the electronic mails stored in the mail folder which has been open.\nAnd, to classify a sent or received electronic mail, the pertinent electronic mail temporarily stored in the mail folder displayed is transferred to a desired mail folder to be appropriately classified. When the entered electronic mail is once stored into the appropriate mail folder, it is shown only when the folder in which it is classified is opened.\nAs described above, the conventional electronic mail cataloging and retrieving system displays in a specified display area a list of header information on the electronic mails stored in the mail folder which is open and attracting attention, and switches to a focused mail folder in order to display a list of header information on the electronic mails stored in another mail folder. Therefore, header in formation on all electronic mails stored in the electronic mail cataloging and retrieving system can not be seen at the same time. When an electronic mail stored in a mail folder different from the mail folder being displayed is wanted to see, the operation of renewing the contents to be displayed in the area for showing a list of header information on electronic mails by switching the focused mail folder is frequently repeated, taking lots of time and labor. Besides, when searching electronic mails based on a meaning which does not correspond to any meaning in the mail folders, search must be done in every folders and that causes the operation of switching folder frequently. Moreover, when reclassifying electronic mails in any mail folders based on a meaning which does not correspond to any meaning in the mail folders, same operation described above is frequently repeated in order to transfer each electronic mail in some folder to other.\nAnd, since header information on electronic mails classified in a mail folder is displayed only when the pertinent mail folder in which they are classified is opened, if a certain electronic mail is erroneously classified into a mail folder, the conventional electronic mail cataloging and retrieving system has a disadvantage that retrieval of this electronic mail involves heavy difficulties.\nIn addition, the conventional electronic mail cataloging and retrieving system has following disadvantages that it takes lots of time and effort to classify electronic mails, when the conventional electronic mail cataloging and retrieving system receive or send mails, it stores at first all these mails in one particular folder generally called as Inbox. Therefore, in order to appropriately classify the electronic mails as described above, the user has to check each electronic mail to classify according to its content, taking lots of time to classify the electronic mails when there are many sent and received electronic mails. In order to remedy such a disadvantage, there have been proposed a method for automatically classifying electronic mails according to classification rule defined by the user, a method for extracting a keyword by automatically analyzing the content of each mail, a method for automatically studying classification rule and the like. Among such conventional technologies, a technology of automatically classifying electronic mails according to user defined classification rules is disclosed in Japanese Patent Laid-Open Publication (Kokai) Nos. Heisei 6-162085, Heisei 2-170642, and Heisei 4-351134. And, a technology of extracting a keyword and automatically studying a classification rule is disclosed in Japanese Patent Laid-Open Publication (Kokai) No. Showa 63-175965.\nBut, the above-described electronic mail classifying technologies have the following disadvantages. Specifically, according to the conventional technology for classifying electronic mails according to a classification rule defined by the user, the classification rule is entered by the user through an input device such as a keyboard. Therefore, the user has to study the forms of classification rules in advance. And, as the application range of electronic mails is expanding these years, the number of rules which have to be entered is increasing, and the rules are becoming complex. Besides, since the electronic mail classification standard of each user is frequently changed, the classification rules have to be changed every time the classification standard is changed. Therefore, the conventional method of entering the classification rules through a keyboard needs increased time and trouble to prepare or change the classification rules.\nAnd, since the conventional technology of automatically studying the classification rules in view of the contents of electronic mails requires a dictionary with enormous amounts of words and background knowledge, it is not suitable to be used with a personal computer because its disk capacity and processing speed are limited. And, since large amounts of electronic mails is required to study the classification rules, it is difficult to immediately correspond to the user's classification standard which is changed in a short period."}
{"text": "The present invention generally relates to an optical disk and a recording/reproducing apparatus therefor, and particularly relates to a tracking system in the recording/reproducing apparatus.\nA pre-groove is formed in a writable optical disk so that an information signal can be recorded efficiently on the optical disk by means of a light while the light beam tracks on the pre-groove. As a system of making a light beam perform tracking on a pre-groove formed in an optical disk, a far field system or a sample tracking servo system using wobbled pits has been conventionally used.\nIn the far field system, however, there has been a disadvantage in that an offset is sometimes generated in a tracking error signal. Further, in the sample tracking servo system, on the other hand, there has been a disadvantage in that a control circuit is complicated and a servo frequency band is limited (for example, Japanese Patent Post-examined Publication Nos. 59-18771 and 63-43812)."}
{"text": "Field of the Invention\nThe present invention relates to a backlight unit and a liquid crystal display device including the same, and more particularly, to a direct-type backlight unit and a liquid crystal display device including the same.\nDiscussion of the Related Art\nPortable electronic devices such as mobile terminals and notebook computers typically include flat panel display devices.\nExample flat panel display devices include a liquid crystal display device, a plasma display panel, a field emission display device, a light emitting diode display device, and an organic light emitting diode display device.\nAmong the flat panel display devices, the liquid crystal display device has advantages such as highly developed mass-production technology, easy driving means, low power consumption, high definition resolution, and large screen size.\nSince the liquid crystals themselves do not emit light, the liquid crystal display device includes a separate light source such as a backlight unit provided with a light-emitting diode (LED) on a rear surface to irradiate light toward the surface of the liquid crystal panel in order to display an identifiable image.\nBacklight units may be categorized into direct-type backlight units and edge-type backlight units. In a direct-type backlight unit, light sources are arranged below a liquid crystal panel. In an edge-type backlight unit, a light source is arranged to face a side of a light guide plate, and a plurality of optical sheets are arranged between the liquid crystal panel and the light guide plate.\nFIG. 1 is a cross-sectional view illustrating a related art liquid crystal display device that includes a direct-type backlight unit. The direct-type backlight unit has a structure where a plurality of optical sheets 5 and a diffuser plate 4 are disposed below a liquid crystal panel 6, and a plurality of light sources 1 are arranged below the diffuser plate 4. Light generated from the light sources 1 is scattered and refracted through the diffuser plate 4 and the optical sheets 5, and then spreads toward an entire surface of the liquid crystal panel 6. The liquid crystal display device also includes a bottom cover 7, a guide panel 8, and a case top 9 to secure the light sources 1, diffuser plate 4, optical sheets 5, and liquid crystal panel 6.\nIf the light sources 1 are arranged below the diffuser plate 4 as shown in FIG. 1, a deviation of light supplied to an edge area of the liquid crystal panel 6 and the other areas is generated, which results in non-uniform display color across the surface of the liquid crystal panel 6.\nFIG. 2 is a cross-sectional view illustrating a related art liquid crystal display device that includes a direct-type backlight unit having a support side 3. As shown in FIG. 2, support sides 3 having a shape tilted with respect to bottom cover 7 are arranged at both sides of the light sources 1. A path of light generated from the light sources 1 is changed by the support side 3 in order to increase light supplied to an edge area of liquid crystal panel 6. The liquid crystal display device also includes a guide panel 8 and a case top 9 to secure the diffuser plate 4, optical sheets 5, and liquid crystal panel 6.\nAs the liquid crystal display device's bezel becomes narrow, the support side 3 is arranged on a rear surface of an active area where an image is displayed on the liquid crystal panel 6 as shown in FIG. 2. This results in non-uniformity of light in area S between the light source 1 and the support side 3, which causes the user to perceive relative darkness in area S.\nFIG. 3 illustrates an image of light provided to a front surface of an area under “A-A” of the direct-type backlight unit of FIG. 2. In FIG. 3, a dark red color indicates that luminance is high as relatively more light is provided, and a faded red color indicates that luminance is low as relatively less light is provided. The amount of light reaching the area S of FIGS. 2 and 3 is less than light reaching the other areas. When an image is displayed by the liquid crystal panel, the image formed in area S is relatively darker than the image formed in other areas.\nAs described above, the related art direct-type liquid crystal display device has a problem that a dark area is generated in the edge area of the liquid crystal panel by non-uniform light supply.\nIt is difficult to provide a large sized screen of 100 inches or more as a single display device. Technology called a VideoWall, MultiVision, or multiple unit display device has been developed.\nThe MultiVision display device is a display device that provides a large sized screen by continuously arranging a plurality of unit display devices in horizontal and vertical directions. The MultiVision display device displays different images on each unit display device or displays one image which is divided on each unit display device.\nParticularly, when multiple images are displayed by a display device, more important images are generally displayed in a center area of the display device and the less important images are displayed towards edge areas of the display device.\nTherefore, if one image is divided for display on unit display devices in the MultiVision display device, the more important image is displayed by the unit display device arranged at the center of the overall display device. As a result, a dark area is generated in the edge area of the unit display device that is arranged at the center.\nIn this way, if one image is divided for display through unit display devices of the MultiVision display device, problems occur in that a dark area may occur in an important image, and a user may perceive quality deterioration more seriously than when the unit display devices each display different images."}
{"text": "The invention is based on a precontrolled 3-way pressure control valve.\nA 3-way pressure control valve of kind is known from the book \"Hydraulik Trainer Band 1\" [Hydraulics Trainer Volume 1], published in 1991 by Mannesmann Rexroth GmbH, pages 236 to 239. In this pressure control valve, the main piston has a piston bead by which the middle control chamber can be sealed off from the first and second control chambers and which is abutted on either side by a piston neck from which control grooves distributed around the circumference extend into the piston bead. The control grooves that begin at a piston neck are located at an axial distance from the control grooves departing from the other piston neck, said distance being approximately equal to the axial extent of the middle control chamber.\nThe pilot valve of the known precontrolled pressure-regulating valve is a direct-controlled pressure-regulating valve with a pressure connection, a tank connection, and a control connection. The control connection is connected with a pressure chamber in front of one end of the main piston by fluid. Control oil from the pressure connection of the main valve flows to the pressure connection of the pilot valve. Depending on the adjustment of a spring, a certain pressure is set in the control connection of the pilot valve and hence in front of one end of the main piston. The middle control chamber of the valve bore of the main valve is connected by a housing channel with a pressure chamber in front of the other end of the main piston. In each case, a pressure becomes established in the middle control chamber and hence in the control connection such that an equilibrium exists at the main piston between the various pressure and spring forces. A precontrolled 3-way pressure control valve in which the pilot valve is supplied with the control oil from the pressure connection of the main valve is referred to as being high-pressure-controlled. The special advantage of this type of precontrol consists in the main piston not tightening even with a large volume flow and a large pressure differential between the pressure connection and the control connection of the main valve, i.e. the connection between the pressure connection and the control connection is interrupted by the main piston and hence the applicability of the valve is limited. It is disadvantageous that the known precontrolled pressure control valve has a tendency toward instability under certain operating conditions, in other words pressure fluctuations occur at the control connection."}
{"text": "Typical gas-fired water heaters are constructed for installation and operation in indoor spaces such as basements, laundry rooms and closets, for example. As a result, the materials selected for water heater components and construction of such water heaters do not lend themselves to installation and operation of water heaters outside of such buildings. Exterior operation subjects water heaters to the elements which can reduce water heater longevity and reduce operating efficiency.\nThere is a growing need for water heaters capable of exterior installation and operation in view of increasing regulation of placement and operation of water heaters in interior spaces. For example, many localities now have regulations concerning the need to supply combustion air from outside the structure instead of the traditional means of supplying combustion air from the interior space itself. A number of water heaters have been developed to address these problems, such as the water heater disclosed in U.S. Pat. No. 5,020,512. That water heater uses a first concentric tube to supply combustion air from outside of the building in which the water heater is located and a second concentric tube to exhaust combustion gases outwardly of the building. Combustion air is channeled from the outer concentric tube into an outer jacket surrounding the water heater, the outer jacket forming a space communicating with the combustion chamber of the water heater.\nSuch water heaters typically require cumbersome concentric or multiple tubing with the necessary associated boring or cutting through the wall of the structure and then mounting vents to the wall. This additional apparatus is expensive and sometimes not easy to install because of the distance to the wall or obstacles between the wall and the water heater. Thus, exterior installation would be a possible alternative. However, such water heaters are constructed from materials suitable for indoor use which does not account for harsh exterior conditions likely to be encountered by an outdoor water heater.\nAnother problem necessitating water heaters for exterior installation and operation is the lack of suitable interior space for water heaters. Many dwellings are constructed without basements and, since space is at a premium, many dwellings seek to minimize or eliminate the space occupied by bulky water heaters. Multiunit dwellings are especially frequently space deficient and can benefit by exterior installation of water heaters. Also, since many dwellings are constructed without masonry chimneys, exhausting the hot flue gases can become a significant problem if the water heater is not designed to reduce flue gas temperatures.\nAlthough many typical water heaters have been installed outside of living spaces, such as in garages, out buildings and the like, attempts to employ such water heaters completely outside have not been successful. The materials of construction of the water heaters lend themselves to premature corrosion of exterior parts. Installation of water heaters outside severely reduces energy efficiencies, especially in northern locations wherein combustion air often is very low in temperature and the entire unit is subjected to very severe low temperatures. Installation of water heaters outside can lead to reduced operating efficiencies because of the cold temperatures or because of excess wind blowing out pilot lights and accumulation of debris, such as leaves and the like around the base of the water heater, thereby reducing a balanced supply of combustion air to the water heater."}
{"text": "Silicates have been used in electrocleaning operations to clean steel, tin, among other surfaces. Electrocleaning is typically employed as a cleaning step prior to an electroplating operation. Using \"Silicates As Cleaners In The Production of Tinplate\" is described by L. J. Brown in February 1966 edition of Plating.\nProcesses for electrolytically forming a protective layer or film by using an anodic method are disclosed by U.S. Pat. No. 3,658,662 (Casson, Jr. et al.), and United Kingdom Patent No. 498,485; both of which are hereby incorporated by reference.\nU.S. Pat. No. 5,352,342 to Riffe, which issued on Oct. 4, 1994 and is entitled \"Method And Apparatus For Preventing Corrosion Of Metal Structures\" that describes using electromotive forces upon a zinc solvent containing paint."}
{"text": "Commonly known wheeled carrying cases include a substantially rigid frame providing side walls of the case, and front and rear major walls that are supported by the frame, at least one of the major walls including a zippered attachment to adjacent ones of the side walls.\nWheeled carrying cases are also well known in the art that include a handle that can be telescopically extended, or retracted into, the body of the case at one of the side walls thereof, and, wheels provided at the opposite side wall of the case that can roll over a ground surface, such as the floor of an airport terminal, or, a passenger aisle of an aircraft.\nWhile this arrangement is admirable for relatively small articles of luggage, it poses a problem in the event that a large and heavily loaded article of luggage is involved.\nWhile such large and heavy articles of luggage easily can be stored in the luggage compartment of an automobile or other form of conveyance, when it is necessary to transport the article of luggage to some other location, a problem arises in providing adequate stability for the article of luggage when being wheeled on a ground surface, this being due to the fact that during wheeling of the article of luggage, the article of luggage must be held in an inclined position. In that position, the center of gravity of the case is displaced laterally of the ground engaging wheels, is located between the handle and the ground engaging wheels.\nThis can constitute a source of discomfort to the person towing the article of luggage, who is then required to exert sufficient force to support the handle, which at the time the article of luggage is positioned in the inclined position relative to the vertical for towing, can exert a considerable downwards force on the handle, that downwards force being resisted by the grasp of the user's hand.\nFurther, depending on the physical height of the person towing the article of luggage, the condition exists that a relatively shorter person will be subjected to a much greater force, due to the greater inclination relative to the vertical of the article of luggage and the fixed position of the handle relative to the article of luggage, than would be a taller person, who conveniently can hold the article of luggage at a lesser angle of inclination relative to the vertical during towing of the article of luggage.\nAs will be appreciated, the wheels on which the article of luggage is towed act as a fulcrum about which the article of luggage angles relative to the vertical when moved from the vertical position to the inclined towing position. Further, the greater the angle of inclination relative to the vertical, the greater becomes the spacing of the center of gravity relative to the towing wheels providing the fulcrum, the center of gravity then being displaced laterally of the vertical and laterally of the axis of the ground engaging wheels.\nOne manner of reducing this problem would be for the user to have packaged all heavy articles at the bottom of the article of luggage and in proximity to the towing wheels, the remainder of the case being occupied by relatively light weight articles. This, however, does not normally occur during the packaging of the article of luggage by the user. If relatively heavy articles are packaged in that end of the article of the luggage remote from the towing wheels, then, the center of gravity of the packaged article of luggage becomes even more displaced relative to the vertical plane including the axes of the towing wheels.\nIf the handle of the case should slip out of the grasp of the user's hand, then, the only possible situation is that the article of luggage, and its contained articles will crash downwardly onto the supporting ground surface, this resulting in possible damage to the articles stored within the article of luggage.\nAnother problem often encountered by travelers is having to deal with multiple articles of luggage. Even if both have wheels, it is very cumbersome to attempt to tow (or push) more than one article of luggage at a time. As well, the ability to stack one article of luggage upon another article of luggage is restricted and also adds significant weight along the extent of the inclined handle of the towed article of luggage.\nFor this reason, despite the generally recognized advantages of wheeled luggage, there are significant unresolved problems in use which the simple addition of wheels and a handle do not solve."}
{"text": "This invention relates to optical position encoders and, more particularly, to arrangements of optical fibers or light guides that function as a readhead for such optical position encoders.\nAs is known in the art, optical encoders include a plurality of adjacent strip-like bands that are subdivided into regions or segments that exhibit first and second reflection coefficients (e.g., transmissive regions and highly reflective regions). When the optical encoder is illuminated along a linear path that traverses the optical encoder bands, light reflected from each band is detected with the two degrees of reflectivity corresponding to the \"0\" and \"1\" binary states. Collectively, the light reflected from the encoder bands thus forms a binary encoded optical signal that is representative of the position of the encoder at the time of illumination.\nOne type of prior art arrangement for illuminating the bands of an optical encoder and, in addition, receiving the reflected light includes a linear array of optical fibers that are arranged and positioned so that the end regions of the fibers are maintained substantially parallel to one another with the end faces of the fibers being in spaced apart parallel juxtaposition with the surface of the optical encoder. In such an arrangement, which commonly is called a reading head or readhead, each band of the optical encoder is spatially and operably associated with one or more specific optical fibers that direct light onto the surface of the associated encoder band and receive light reflected from the associated encoder band. In some arrangements, the associated optical fibers both transmit light to the surface of the optical encoder and receive reflected light. In other arrangements, separate transmitting and receiving fibers are used.\nIn operation, a pulse of light is transmitted along those optical fibers that are employed to transmit light to the associated encoder bands. Light, launched from the ends of the optical fibers impinges on the surface of the optical encoder with light that is reflected from each individual band being received either by the transmitting fibers or by dedicated receiving fibers. The reflected light then propagates along the associated receiving fibers for subsequent reception and processing to form a binary encoded optical or electrical signal that is representative of the encoder position relative to the fixed or reference position of the readhead.\nFor example, a rotary optical encoder for producing an N-bit Gray Code signal representative of the angular position of a shaft or other rotatable object generally is formed as a disk that includes N circumferential bands concentrically surrounding the axis of rotation. Each circumferential band is subdivided into interspersed highly reflective and highly transmissive equal angular radial segments with the two innermost circumferential bands each including a single reflective (and transmissive) segment that is subtended by an angle of .pi. radians (180.degree.). In this arrangement, the number of reflective (and transmissive) segments in the N-2 circumferential bands located outwardly of the two innermost bands sequentially increases in accordance with the sequence 2.sup.i-2 (i=3, 4, . . ., N); where i identifies the circumferential band (as counted from the center of the optical encoder). Thus, each reflective (and transmissive) segment of the outmost N-2 circumferential bands is subtended by an angle of .pi./2.sup.i-2. In addition, the reflective (and transmissive) regions of the two innermost circumferential bands are positioned relative to one another so that the reflective (and transmissive) segments overlap one another by .pi./2 (90.degree.) and each reflective segment of the outermost N-2 circumferential bands overlaps a contiguous reflective segment of the nextmost inner circumferential band by an angle of .pi./2.sup.i-1. When a suitably configured readhead is mounted at a fixed position that extends radially along the surface of such an encoder disk, an N-bit binary signal is formed by light that is reflected from the encoder bands to uniquely identify the radial position (angle of rotation) of the optical encoder. Since an encoder of the described type provides 2.sup.N binary encoded N-bit binary signals, an angular resolution of .pi./2.sup.N-1 radians is attained.\nAlthough prior art readheads of the above-described configuration perform satisfactorily in many applications, disadvantages and drawbacks are encountered in situations that require optical encoders of relatively small size with an accompanying need for precise determination of angular position (high resolution). In particular, such applications require a relatively large number of small diameter optical fibers which must be mounted in precise relationship both with respect to one another and with respect to the surface of the optical encoder. Because of the small fiber size and positional constraints, small high resolution optical readheads are not well suited to automated large scale fiberoptics manufacturing techniques and are otherwise difficult to fabricate. Thus, in the prior art, such devices have been both costly and, in many cases, not as reliable as desired.\nMoreover, when small high resolution optical readheads are fabricated with prior art fabrication techniques using discrete optical fibers, it often is difficult to achieve desired performance specifications. For example, manufacturing tolerances associated with the assembly of individual fibers may not closely align groups of fibers with the associated bands of the optical encoder. When proper alignment is not achieved, light that is launched from the end of the fibers may impinge upon an adjacent band of the optical encoder and/or light that is received by the optical fibers may have been reflected from an adjacent band of the optical encoder. If this occurs, an undesirable amount of \"cross-talk\" may be encountered (i.e., the optical signals provided to and received from the individual bands of the optical encoder may not be properly isolated from one another). In addition, using prior art discrete fiber fabrication techniques, it may be difficult to obtain low insertion loss and desired efficiency."}
{"text": "1. Field\nEmbodiments relate to a power control device and an image forming apparatus including the same. More specifically, embodiments relate to a power control device configured to reduce standby power consumption, and an image forming apparatus including the same.\n2. Description of the Related Art\nFrom the point of view of environmental protection, restriction on standby power consumption of electrical appliances for energy-saving is being strongly imposed in Americas (EPA1.2) and Europe (ErP step 2) as well as in Korea. In order to cope with the restriction, various efforts are under way to reduce the standby power consumption of electrical appliances.\nFor example, an electromagnetic interference (EMI) filter for removing noise is included in an input of a power supply (e.g. a switch-mode power supply), and if electrical charges of an X-capacitor included in the EMI filter flow to a metal terminal of a plug during a plug-off state, a safety problem may occur. Accordingly, in order to solve the problem, a discharge resistor configured to discharge charges filled in the X-capacitor may be used. However, there is power consumption in a plug-on state due to the discharge resistor, which is against the effort for reducing the standby power consumption of electrical appliances. Therefore, measures to solve the problem are needed."}
{"text": "This invention relates to a microprocessor, and more particularly to a microprocessor and associated control method capable of dealing with both large-scale intricate program processing and fast program processing which is intended for real-time operation.\nA microprogram-based processor is disclosed in Japanese Patent Unexamined Publication No. JP-A-56-108149, for example. This processor is based on a control scheme in which microinstruction routines common to many different macroinstructions are stored in a nano-instruction memory and used commonly for the execution of macroinstructions.\nA microcomputer capable of running multiple independent programs is disclosed in the co-pending U.S. patent application Ser. No. 07/414,315 filed on Sep. 29, 1989 for example. This microcomputer includes a main processor and a subordinate processor. The main processor is a usual processor such as the one disclosed in the above patent publication JP-A-56-108149, for example, and the sub processor is a processor intended for real-time processing with the ability of running multiple independent programs quickly.\nThe processor of the above patent publication 56-108149 is designed to run multiple independent programs in such a manner that each program is created in the form of sub programs which are linked by a main program. However, because of sequential operations of sub programs, it is difficult for this processor to deal with fast real-time processing in which the processor needs to respond to input signals promptly.\nThe microcomputer disclosed in the co-pending application Ser. No. 07/414,315 is operative for both the usual program processing and fast real-time processing. However, this microcomputer inevitably has a large-scale circuit arrangement due to the two processors included in it."}
{"text": "1. Field of Invention\nThe present invention relates to a method for drawing an optical fiber having reduced polarization mode dispersion (PMD). The present invention also relates to the single mode optical fiber made in accordance with the presently disclosed method.\nSingle mode optical fiber commonly used in communication systems is not purely single mode. Rather, two modes, with perpendicular polarizations, exist in single mode fiber. These two polarizations form an orthogonal basis set. Accordingly, any configuration of light that propagates through a single mode fiber can be represented by a linear superposition of these two modes.\nIf the fiber has perfect circular symmetry in both geometry and internal applied stress, the two polarization modes propagate with the same group velocity. They have no time delay difference after traveling the same distance in the fiber. But a practical fiber does not have perfect circular symmetry. Imperfections such as geometric deformation and stress asymmetry makes the two polarization modes to propagate with different velocities (a function of propagation constant). The difference between the propagation constants is termed as birefringence. The differential time delay between the two polarization modes is called PMD (polarization mode dispersion). This limits the high bit rate transmission in the communication system.\n2. Description of Prior Art\nAccording to prior art of spinning the fiber mentioned in U.S. Pat. No. 5,418,881; the fiber is subjected to a sinusoidal spin function, which varies with the spatial distance of the fiber drawn. However, this patent does not teach about the discrete spin rates and a trapezoidal wave function of the spin rate with the spatial distance of the fiber.\nIn another published PCT application bearing No. WO/97/26221, the authors have identified the disadvantage of above-cited U.S. Pat. No. 5,418,881. The major disadvantage identified is being in the difference of beat lengths along the different lengths of the fiber. A particular combination of spin amplitude and spin frequency will lower the PMD effectively for a particular beat length. However, this published PCT application WO/97/26221 does not teach the manner for spinning the fiber, particularly after coating and twisting the fiber so as to achieve effective lowering of the PMD of the fiber, particularly cabled fiber. The author of this published PCT application has not visualized that the use of combination of discrete spin rates of different values, which repeats itself in a periodic fashion can effectively lower the PMD of the cabled fiber.\nThe present invention now makes a disclosure of a method by which it is possible to achieve the reduced PMD in optical fiber as low as less than about 0.1 ps/√km and the post cabling PMD to be as low as less than about 0.3 ps/√km when the spinning is performed on a fiber, particularly on a drawn, coated and twisted fiber by employing the spin function, which is a combination of four discrete spin rates which repeat itself in a periodic fashion after a certain length of the fiber drawn and the spin rate of each spin function varies in a trapezoidal manner along the drawn length of fiber. The fibers drawn from performs in accordance with the method of the present invention have ovality at max about 3%."}
{"text": "This invention concerns a duet-sound generating method for use in an electronic musical instrument.\nGenerally, a duet-sound refers to an ensemble of notes. In other words, it may be a sound produced from incorporating a current melody note being played with a machine generated note or notes. To generate such a duet-sound note interval, as shown in FIG. 1, a conventional electronic musical instrument includes: a microcomputer 11 that controls related sound-source control data and generation thereof by inputting various control-switching signals and melody keyboard key switching signals; an auto-accompaniment switch 13, a duet switch 15, a filter selection switch 19, and a melody keyboard 17 that generates depressed key note signals. A sound-source generation integrated circuit (referred to as a sound-generator IC hereinafter) 21 generates related sound-source and chord sounds according to the chord and the sound-source control data from the microcomputer 11 and outputs them to the output terminals OUT1 and OUT2; a chord and bass filter 41 that, enabled under the control of the microcomputer, filters the output of the sound-generator IC's chord and bass; a rhythm generator a 43 generating rhythm signal under the microcomputer's control; a plurality of filters 45, 47 and 49 for filtering the sound-source output of the sound-generator IC 21; a pre-amplifier 51 that pre-amplifies the outputs of the filters; pre-amplifiers 53, 55 that pre-amplify the outputs of the chord and bass filter 41 and the rhythm generator 43; and a main amplifier 57 that amplifies the output of the pre-amplifiers 51, 53, 55 and outputs them to the speaker 59.\nFIGS. 2a and 2b are flow chart for a conventional duet-sound generation, that is a program for use in the microcomputer 11 of FIG. 1. FIG. 2a represents the auto-accompaniment part and FIG. 2(b) is the non-auto-accompaniment. In auto-accompaniment mode's case, the duet sound is programmed to be produced when a melody note on the melody keyboard 17 of FIG. 1 is depressed and the duet switch 15 is turned `on`. The conventional duet-sound generation procedure is explained in the following section with reference to FIGS. 1 and 2. When the Filter Selection switch 19 is switched to Filter 1, e.g., a piano-voice selection signal to the microcomputer 11, the microcomputer recognizes the Filter 1 selection condition (e.g., paino sound) performs key scanning and sends a logic high signal to the output terminal S1 to enable Filter 1, filter 45. If the auto-accompaniment switch 13 becomes \"on\" under such a condition, the microcomputer 11 recognizes the change by key-scanning activity and determines the current mode is auto-accompaniment on. It then sends chord data for the generation of chord sounds to the sound-generator IC 21 and outputs the chord and bass-filter enable signal and rhythm-enable \"high\" signal at terminal S8 to enable them. During this time, the sound-generator IC 21 produces and sends the related chord sound to the output terminal OUT2, which causes the chord and bass filter 41 to output the filtered chord sound. The enabled rhythm generator 43 produces a certain rhythm. The chord and rhythm signals generated above, go through respective pre-amplifiers 53 and 55 to be amplified to a certain level, mixed in the main amplifier 57, and finally sent to speaker 59. During operation, pressing a key on the melody keyboard 17 prompts the microcomputer to check if the duet switch 15 is on. If the duet switch 15 is \"off\", a channel is assigned to a key to produce a sound. But, with Duet Switch 15 on, the melody keyboard 17 is scanned for a root note for the output chord. As in the logic of FIG. 2a, if a melody keyboard note is depressed without the root note present data for the root note and the melody note are allocated to two channels to be inputted to the sound-generation IC 21. At the moment, the sound-generation IC 21 generates the sound source according to the microcomputer's data and sends it to the output terminal OUT1. The sound-generation IC's sequentially outputting the data allocated to the two channels by the microcomputer 11, thereby generating the duet sound. For example, if a chord currently being outputted is a C chord, a \"Do\" sound that is equivalent to the root note of the chord is allocated to one of the two channels and a \"Mi\" sound is allocated to the other one, thereby being outputted to the sound-generator IC's output terminal OUT1. The duet sound outputted from the sound-generator IC 21 is inputted to the enabled filter 45 to be then processed into a duet sound using Filter 1 upon for the sound of the output. On the other hand, if the result of the above checking operation shows that any melody keyboard note equivalent to the root note of the outputting chord is depressed, then the outputting chord is checked to see if it is major key or a minor key. Thereafter, data for harmonious the major or minor key responding note interval with respect to the depressed key, is sent to the sound-generator IC 21. The sound-generation IC 21, at this moment, produces the sound source corresponding to the above data and sends it to the filter 1, 45 to generate the duet sound.\nThe duet sound produced by the instrument of FIG. 1 according to the flow chart of FIGS. 2a and b, however, provides only a mono-tone single instrument voice sound, incapable of producing duet-sounds with various instrument voices. In other words, the method may accomplish the duet sound only with a sound of a single instrument."}
{"text": "1. Field of the Invention\nThe present invention generally relates to vine crop harvesting apparatus and more particularly concerns an improved oscillatory fruit-vine separator adapted for use with a vine crop harvester.\n2. Description of the Prior Art\nU.S. Pat. No. 3,071,196 to Scheidenhelm discloses a tomato harvester having an oscillatory fruit-vine separator which includes a plurality of endless parallel conveyor chains that are transversely spaced to permit the tomatoes to drop therebetween to an underlying delivery conveyor. The chains are supported by a pair of parallelogram frameworks and are continuously driven while a crankshaft drive mechanism is provided for oscillating the parallelogram frameworks, and thus the conveyor chains. The conveyor chains are provided with outwardly extending fingers which become entangled in the vines to both propel and transmit shaking forces to the vines.\nA problem with the oscillating separators of the aforedescribed type is that the vines tend to become so densely matted adjacent to the conveyor chains that some tomatoes shaken from the vines become entrapped on top of the mat of vines and are thereby prevented from penetrating the vines and falling between the conveyor chains onto the underlying delivery conveyor. It will be appreciated that such entrapped tomatoes will be discharged from the separator onto the field, thereby reducing the yield of the harvester from what it should be.\nU.S. Pat. No. 3,618,617 to Gates et al discloses an oscillatory fruit-vine separator having a plurality of spaced endless belt conveyors mounted on a frame which is supported at its front end by a rockshaft and at its rear end by crank means. The rockshaft provides a fore-and-aft movement to the front end of the frame, and the crank means drives the rear end of the frame in a generally circular motion involving both vertical and fore-and-aft components. To impart a rotary movement to the vines as they are conveyed by the belt conveyors, the Gates et al separator further includes a plurality of tines which are mounted so as to project downwardly from above the separator into upper portions of the vines carried on the belt conveyors. The rods are mounted so that they do not oscillate with the frame carrying the belt conveyors. The patent states that the tines act to impale a vine and slow the progress of part of it, while the rest of the vine continues to be pulled along by the belts of the separator. The rods thereby rotate the vines as they are conveyed."}
{"text": "During the installation and contruction of an offshore underwater pipeline, pneumatic pressure may be applied internally to the pipeline to prevent it from collapsing or buckling from external pressure. Especially in deep water offshore pipeline construction internal pressure is applied to prevent any radial collapse of the line due to the extremely high pressure exerted against it by the surrounding water. In order to limit the weight of the normally unsupported pipeline between the laybarge and the ocean floor, the pressure generally will be applied pneumatically. Sealing devices such as pigs are required internally of the pipeline and usually near the barge in order to maintain the various sections of the assembled pipeline under an internal pressure. In the event that a sealing device malfunctions and does not hold in place the pig normally would be expelled out of the pipeline onto the laybarge like a cannonball because of the potential energy of the pressurized pipeline. This possibility presents a very serious and dangerous situation which the present invention precludes from happening.\nThe specific prior art references considered with respect to the present invention include the following U.S. Pat. Nos.: 2,607,370 3,197,195 2,615,413 3,381,714 2,672,161 3,495,626 2,747,690 3,503,424 2,807,937 3,525,226 2,823,634 3,561,320 2,974,685 3,561,490 3,593,749\nThe foregoing prior art relate to conventional pipeline pigs and plugs and generally disclose the known expedient of providing actuatable means for gripping engagement with the interior wall of the pipe upon occurrence of a predetermined condition. None of the prior art, however, discloses the construction and arrangement of the subject invention per se or in combination with the movable pig which actuates the safety stop apparatus, or the functional interrelationship between the safety stop apparatus which is operable when the sealing pig comes in contact therewith."}
{"text": "Operation of a four-stroke internal combustion engine usually includes an intake stroke, a compression stroke, a power stroke and an expansion stroke. In normal operation, during the intake stroke of the internal combustion engine, an intake valve is opened as a piston is moved downward through a cylinder. At low speed or load engine conditions, a swirl effect or rotation of a fluid mixture around an axis of the cylinder may be used to promote engine performance.\nAlternatively, at high speed or load engine conditions, a tumble effect for a fluid generally represents the fluid mixture spinning around the cylinder along an axis parallel to the crankshaft. Effective fuel economy and performance of an engine may be affected if the tumble effect or swirl effect is not managed properly or applied at an appropriate operation cycle of the engine."}
{"text": "The modification of natural or synthetic polymers with enzymes is an environmentally friendly alternative to classical chemical modification reactions which generally require harsh reaction conditions. Despite the advantage of using enzymes in functionalization reactions, i.e. milder reaction conditions and highly specific transformations, only a few examples have been reported in the literature involving either natural or synthetic polymers.\nMushroom tyrosinase has been used to introduce phenolic functionalities into chitosan. The synthetic pathway involved enzyme-catalyzed oxidation of phenol to an o-quinone which dissociated from the enzyme and freely diffused to the nucleophilic amine site of the chitosan polymer. In addition to the incorporation of quinone into the chitosan, subsequent polymerization of quinone into oligomeric phenols was also observed as determined by UV-Vis spectroscopy. The effectiveness of the reactions was not discussed. Polysaccharides have been functionalized with fatty acids by transesterification using ion-paired subtilisin Carlsberg protease in organic solvents. It has been observed that the enzyme selectively acylated the primary hydroxyl groups on polysaccharides and that the degrees of substitution per glucose moiety in amylose and β-cyclodextrin were 0.185 and 0.250, respectively. Lactose was attached to the hydroxyl groups of hydroxyethylcellulose (HEC) by transglycosylation in sodium acetate buffer using Aspergillus oryzae. Although the number of sites available in each HEC unit is 3, the maximum degree of substitution obtained by this method was only 0.033. The same natural polymer, HEC, has recently been modified by enzymatic transesterification with both vinyl propionate and vinyl acrylate in the presence of subtilisin Carlsberg in anhydrous pyridine. These biotransformations of cellulose were promising; however, low conversions limited their viability. It has been reported that the phosphorylation of cotton cellulose using bakers' yeast hexokinase as the enzyme and ATP as the phosphoryl donor. The phosphorylation of 0.03% of the glycopyranose units in the cellulose resulted in improved dyeability and flame resistance. All of the examples listed here were hindered by low conversion.\nEven fewer examples have been reported for synthetic polymer functionalization using enzymatic catalysis, which are also characterized by low conversion. For instance, 16% of the pendant nitrile groups of polyacrylonitrile fibers were converted to the corresponding amides by a nitrile hydratase enzyme, resulting in a significant increase in the hydrophilicity of the fiber surface. Mushroom tyrosinase has been utilized to graft poly(4-hydroxystyrene) (PHS) onto chitosan. In this approach, first 1-2% of the phenolic moieties of PHS were enzymatically oxidized and then the resulting polymer was reacted with the amine groups of chitosan. Lipase-catalyzed acylation of poly[N-(2-hydroxypropyl)-11-methacryloylaminoundecanamide-co-styrene] and the corresponding monomer with vinyl acetate, phenyl acetate, 4-fluorophenyl acetate and phenyl stearate has been reported. 1H NMR results revealed that the reactivity of the monomer was higher than that of the copolymer and that the copolymer could be acylated with up to 40% conversion. In addition, it was found that the reactivity of copolymer was dependent on copolymer composition which indicated the effect of steric hindrance and hydrophobicity on reaction kinetics. The synthesis of organosilicon carbohydrate macromers by Candida antarctica lipase B catalyzed esterification has been reported. Diacid-endblocked siloxanes were reacted with α,β-ethyl glucoside under vacuum in bulk. Esterification occurred with high regioselectivity (>98%) at the primary hydroxyl (C6) of the glucoside, but electrospray ionization mass spectrometry (ESI MS) showed the presence of a mixture of mono- and diesters.\nThe preparation of polyolefin-based telechelic polymers, such as methacrylate-terminated polyisobutylenes (PIB-MAs), has been shown in the prior art by a variety of strategies. In one example, PIB-Clt prepared by the inifer technique was first converted to the corresponding exo-olefin by dehydrochlorination with t-BuOK, followed by hydroboration/oxidation to yield PIB-OH which was subsequently converted to PIB-MA by acylation with methacryloyl chloride. In another process, PIB-OH was synthesized by hydroboration/oxidation of allyl-terminated PIB, prepared by end-quenching of living PIB+ with allyltrimethylsilane, and the latter was converted to PIB-MA with methacryloyl chloride. A further process describes how PIB-MA was synthesized by nucleophilic substitution of PIB-CH2—CH2—CH2—Br (prepared by anti-Markovnikov hydrobromination of allyl-terminated PIB) with sodium methacrylate. This same method was also applied to Glissopal®2300, a commercially available PIB with olefinic end groups. Quantitative syntheses of primary hydroxyl terminated PIB and Glissopal-OH has also been described.\nNotwithstanding the state of the art as described herein, there is a need for further improvements in the preparation of functionalized polymers (monofunctional or difunctional (telechelic) polymers) via enzymatic catalysis which is an environmentally friendly method when compared with toxic acylating agents such as methacryloyl chloride, and catalysts; and is useful in biomaterials prepared from these functionalized polymers."}
{"text": "The majority of cost-effective digital video cameras use a single sensor with a color filler array (CFA) to capture visual scenes in color. At each pixel, only one of the three primary colors is sampled and the missing color components need to be interpolated. The color interpolation process is usually referred to as color demosaicking (CDM). For video sequences, the temporal correlation in adjacent frames can also be exploited to improve the performance of CDM. The majority of CDM algorithms assume that the CFA data is noise-free. This assumption, however, is unrealistic and the presence of noise in the CFA data will cause serious CDM artifacts that are difficult to be removed by a subsequent denoising process. Many available denoising algorithms are designed for monochromatic/full-color images/videos and they are not directly applicable to the CFA data. Therefore, it is highly desirable to provide new algorithms for full-color video reconstruction from noisy CFA sequences.\nA commonly used strategy is to demosaick the CFA video first and then later denoise the demosaicked video. However, the noise-caused color artifacts generated in the CDM process can be difficult to remove in the following denoising process. Some joint demosaicking and denoising approaches have been proposed. In a joint demosaicking-denoising algorithm, both demosaicking and denoising are treated as an estimation problem and solved by the total least square technique. One computationally more efficient joint demosaicking-denoising scheme first performs demosaicking-denoising on the green channel, and then uses the restored green channel to estimate the noise statistics to restore the red and blue channels. Inspired by the directional linear minimum mean square-error estimation based CDM scheme, a nonlinear and spatially adaptive filler has been proposed by using local polynomial approximation for CDM. This scheme is adapted to noisy CFA inputs for joint demosaicking-denoising.\nIn addition to the two abovementioned strategies, another approach to reproduce full-color images from noisy CFA data is to perform denoising before demosaicking. However, due to the mosaic structure of CFA, existing monochromatic image/video denoising methods cannot be directly applied to the CFA data. The CFA image can be divided into several monochromatic sub-images using the approaches known from the CFA image compression literature. However, such approaches do not exploit the inter-channel correlation which is useful to reduce noise. A wavelet based CFA image denoising method has been proposed. A principal component analysis (PCA) is proposed based CFA image denoising scheme. A CFA block, which contains color samples from different channels, is used to analyze the local CFA image structure. By transforming the noisy signal into the PCA domain, the signal energy can be better clustered and the noise can be more effectively removed. The PCA-based CFA image denoising algorithm effectively exploits the spatial and spectral correlation in each frame. This method was later extended to gray level images."}
{"text": "1. Field of the Invention\nThe present disclosure relates to a phosphor, a method for manufacturing the phosphor, and a light-emitting device using the phosphor.\n2. Description of Related Art\nA light-emitting diode (LED) light-emitting device mainly comprises a combination of a phosphor and a LED chip serving as an excitation light source, and can realize various colors of emission light according to the combination.\nSpecifically, a white LED light-emitting device, which radiates white emission light, comprises a phosphor in combination with a LED chip that emits light in the blue region. For example, the device comprises a phosphor mixture and a LED chip giving off blue light. In practice, a phosphor emitting yellow light, which is complementary to blue light, is often adopted as the phosphor to produce a pseudo-white LED light-emitting device. Further, a three-wavelength white LED device is also developed which comprises a LED chip emitting blue light, a green- or yellow-light emitting phosphor and a red-light emitting phosphor. As an example of the red-light emitting phosphor used in that light-emitting device, K2SiF6:Mn phosphor is known."}
{"text": "As the feature sizes in complementary metal-oxide-semiconductor (CMOS) technology continue to shrink, it becomes increasingly difficult to fabricate the metal interconnects using conventional processing techniques. For example, using a damascene process to fill trenches with copper often results in undesirable effects including poor liner/seed coverage on the trench walls, pinch off at the trench mouth, and reentrant reactive ion etch (RIE) profiles. In addition, the increasing ratio of the liner to copper, copper grain growth, and copper grain scattering phenomena result in increased copper resistivity, which makes the copper less effective as an interconnect material."}
{"text": "Optical fibers are used in a variety of telecommunications applications, including high-capacity data transmission for networks, data centers and server clusters. Conventional high-capacity optical-fiber links typically include cables (e.g., ribbon cables) that include conventional single-core multimode optical fibers, with each fiber carrying a single data channel. Unfortunately, such cables are limited in their connection density and tend to be bulky.\nMulticore optical fibers (“multicore fibers”) have been developed in part to provide higher data capacity. A multicore fiber has multiple cores embedded in a single cladding, and these cores tend to be single-mode cores with diameters less than 10 microns. Multicore fibers thus promise higher density connections as well as greater data capacity than do conventional fibers. However, multicore fibers tend to be difficult to use in high-capacity data transmission applications since the cores are fixed in place and each core needs to be aligned with high accuracy to a corresponding optical transmitter or optical receiver."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus, an image forming method, and a computer-readable recording medium.\n2. Description of the Related Art\nThere are conventionally-known image forming apparatuses that form images on a recording medium, such as sheet, by using colorless and transparent toner in addition to colored toner of CMYK, or the like. Furthermore, in order to utilize the above-described image forming apparatuses, there is a known technology for achieving visual surface effects by attaching control data to document data, such as “control code for providing predetermined glossy texture to the surface of an image although it is hard to be recognized as a pattern” or “control code for specifying a pattern that may be definitely recognized as a pattern”.\nFor example, there is a disclosed printing control device that applies a visual effect to a recording medium (see Japanese Patent Application Laid-open No. 2012-083736). Japanese Patent Application Laid-open No. 2012-083736 discloses the configuration of the printing control device that processes the gloss-control plate data that specifies the gloss control value for determining the type of surface effect to be applied and its area on the recording medium on a pixel by pixel basis.\nHowever, according to the above-described conventional technology for applying visual surface effects, the control code, processed by a raster image processor (RIP), is replaced with a pattern. Therefore, if images are recorded after the document data is reduced in size, e.g., images are combined before they are formed on a recording medium, the pattern of the transparent toner is used without being reduced in size; therefore, there is a problem in that it is difficult to achieve the desired visual surface effect.\nIn view of the conventional problem, there is a need to provide an image forming apparatus, an image forming method, and a computer-readable recording medium having a computer program that make it possible to achieve the desired visual surface effect even in a case where an image is recorded on a recording medium after the document data is reduced or enlarged in size."}
{"text": "1. Field of Technology\nThe present disclosure relates to equipment and techniques for melting metals and metallic alloys (hereinafter “alloys”). The present disclosure more specifically relates to equipment and techniques utilizing electrons to melt metals and alloys and/or the materials while in a molten state.\n2. Description of the Background of the Technology\nThe alloy melting process involves preparing a charge of suitable materials and then melting the charge. The molten charge or “melt” may then be refined and/or treated to modify melt chemistry, remove undesirable components from the melt, and/or affect the microstructure of articles cast from the melt. Melting furnaces are powered by either electricity or the combustion of fossil fuels, and selection of a suitable apparatus is largely influenced by the relative costs and applicable environmental regulations, as well as by the identity of the material being prepared. A variety of melting techniques and apparatus are available today. General classes of melting techniques include, for example, induction melting (including vacuum induction melting), arc melting (including vacuum arc skull melting), crucible melting, and electron beam melting.\nElectron beam melting typically involves utilizing thermo-ionic electron beam guns to generate high energy substantially linear streams of electrons which are used to heat the target materials. Thermo-ionic electron beam guns operate by passing current to a filament, thereby heating the filament to high temperature and “boiling” electrons away from the filament. The electrons generated from the filament are then focused and accelerated toward the target in the form of a very narrow, substantially linear electron beam. A type of ion plasma electron beam gun also has been used for preparing alloy melts. Specifically, a “glow discharge” electron beam gun described in V. A. Chernov, “Powerful High-Voltage Glow Discharge Electron Gun and Power Unit on Its Base”, 1994 Intern. Conf. on Electron Beam Melting (Reno, Nev.), pp. 259-267, has been incorporated in certain melting furnaces available from Antares, Kiev, Ukraine. Such devices operate by producing a cold plasma including cations which bombard a cathode and produce electrons that are focused to form a substantially linear electron beam.\nThe substantially linear electron beams produced by the foregoing types of electron beam guns are directed into the evacuated melting chamber of an electron beam melting furnace and impinged on the materials to be melted and/or maintained in a molten state. The conduction of electrons through the electrically conductive materials quickly heats them to a temperature in excess of the particular melting temperature. Given the high energy of the substantially linear electron beams, which can be, for example, about 100 kW/cm2, linear electron beam guns are very high temperature heat sources and are readily able to exceed the melting and, in some cases, the vaporization temperatures of the materials on which the substantially linear beams impinge. Using magnetic deflection or similar directional means, the substantially linear electron beams are rastered at high frequency across the target materials within the melting chamber, allowing the beam to be directed across a wide area and across targets having multiple and complex shapes.\nBecause electron beam melting is a surface heating method, it typically produces only a shallow molten pool, which can be advantageous in terms of limiting porosity and segregation in the cast ingot. Because the superheated metal pool produced by the electron beam is disposed within the high vacuum environment of the furnace melting chamber, the technique also beneficially tends to degas the molten material. Also, undesirable metallic and non-metallic constituents within the alloy having relatively high vapor pressures can be selectively evaporated in the melting chamber, thereby improving alloy purity. On the other hand, one must account for the evaporation of desirable constituents produced by the highly-focused substantially linear electron beam. Undesirable evaporation must be factored into production and can significantly complicate alloy production when using electron beam melting furnaces.\nVarious melting and refining methods involve the electron beam melting of feed stocks using thermo-ionic electron guns. Drip melting is a classic method used in thermo-ionic electron beam gun melting furnace for processing refractory metals such as, for example, tantalum and niobium. Raw material in the form of a bar is typically fed into the furnace chamber and a linear electron beam focused on the bar drip-melts the material directly into a static or withdrawal mold. When casting in a withdrawal mold, the liquid pool level is maintained on the top of the growing ingot by withdrawing the ingot bottom. The feed material is refined as a result of the degassing and selective evaporation phenomena described above.\nThe electron beam cold hearth melting technique is commonly used in the processing and recycling of reactive metals and alloys. The feedstock is drip melted by impinging a substantially linear electron beam on an end of a feedstock bar. The melted feedstock drips into an end region of a water-cooled copper hearth, forming a protective skull. As the molten material collects in the hearth, it overflows and falls by gravity into a withdrawal mold or other casting device. During the molten material's dwell time within the hearth, substantially linear electron beams are quickly rastered across the surface of the material, retaining it in a molten form. This also has the effects of degassing and refining the molten material through evaporation of high vapor pressure components. The hearth also may be sized to promote gravity separation between high-density and low-density solid inclusions, in which case oxide and other relatively low-density inclusions remain in the molten metal for a time sufficient to allow dissolution while high density particles sink to the bottom and become trapped in the skull.\nGiven the various benefits of conventional electron beam melting techniques, it would be advantageous to further improve this technology."}
{"text": "1. Field of the Invention\nThis invention relates to helmet mounted displays. In this specification the term \"helmet\" means any article worn by an observer on his/her head and includes protective helmets such as those worn by aircraft pilots or vehicle drivers and simple instrument supporting head bands such as those worn by surgeons and other medical workers. The term \"helmet mounted display\" shall mean apparatus attached to, located on, or forming part of such a helmet and for presenting images or information to the wearer of the helmet.\n2. Description of the Related Art\nFor combat aircraft use, flight testing and medical examination have confirmed that binocular display to the pilot both aids performance and also reduces brain fatigue.\nAt present, binocular presentation is provided by either having:\n(i) two sources of display or PA1 (ii) one source with optical train containing a split prism.\nThe adoption of solution (i) doubles the number of devices on the helmet providing the display and hence increases the helmet mass.\nThe adoption of solution (ii) halves the level of brightness of the display as inherently the split prism has a 50% transmission characteristic.\nConsequently both the existing techniques have operational disadvantages. Firstly, excess mass and the adverse physiological effect under g forces; and secondly, low brightness and hence loss of display under high ambient light conditions.\nThere is alternatively often a requirement to provide the pilot with information f rom two similar or dissimilar data sources simultaneously. Such sources may be for example cathode ray tube, dot matrix, or conventional optical displays. Currently the pilot may view one source only at a time."}
{"text": "Embodiments of the invention relate to neuromorphic and synaptronic computation and in particular, mapping neural dynamics of a neural model on to a lookup table.\nNeuromorphic and synaptronic computation, also referred to as artificial neural networks, are computational systems that permit electronic systems to essentially function in a manner analogous to that of biological brains. Neuromorphic and synaptronic computation do not generally utilize the traditional digital model of manipulating 0s and 1s. Instead, neuromorphic and synaptronic computation create connections between processing elements that are roughly functionally equivalent to neurons of a biological brain. Neuromorphic and synaptronic computation may comprise various electronic circuits that are modeled on biological neurons.\nIn biological systems, the point of contact between an axon of a neuron and a dendrite on another neuron is called a synapse, and with respect to the synapse, the two neurons are respectively called pre-synaptic and post-synaptic. The essence of our individual experiences is stored in conductance of the synapses. The synaptic conductance changes with time as a function of the relative spike times of pre-synaptic and post-synaptic neurons, as per spike-timing dependent plasticity (STDP). The STDP rule increases the conductance of a synapse if its post-synaptic neuron fires after its pre-synaptic neuron fires, and decreases the conductance of a synapse if the order of the two firings is reversed."}
{"text": "This invention relates to a device adapted to be connected to an electrical appliance, such as a clock radio, to visually and audibly give an alarm in the event of a power failure.\nClock radios are commonly used to sound an alarm to awaken a sleeping person at a particular predetermined time. In the event of a power failure unbeknown to the sleeping person, the clock will be turned off during the period of failure and restored after the failure has been corrected. The clock will therefore indicate the incorrect time and the alarm to awaken the sleeping person will be sounded at a time later than desired by an amount equivalent to the period of the failure. This invention provides an alarm device which will provide an audible signal at the time of the power failure to awaken the sleeping person and alert him to the fact of the failure, so the clock can be reset to compensate for the time period of the failure. In the event the audible alarm is not heard, a visual alarm is simultaneously activated to indicate to the sleeping person that a power failure has occurred when he awakens."}
{"text": "1. The Field of the Invention\nThe present invention relates to sodium hypochlorite and other hypohalite or hypohalous acid compositions (e.g., aqueous solutions) that include a cationic polymer, as well as methods of making and using such compositions.\n2. Description of Related Art\nConsumers recognize sodium hypochlorite as a highly effective cleaning, bleaching and sanitizing agent that can be widely used in cleaning and sanitizing various hard and soft surfaces, in laundry care, etc. Often, sodium hypochlorite solutions in a relatively concentrated form are delivered to the consumer, who then dilutes them. A wide variety of automated washing machines which can control the addition of sodium hypochlorite to laundry washwater are also available. When used in washing machines, the sodium hypochlorite solution will typically come in contact with the detergent being used in the laundering process. Complex interactions between sodium hypochlorite and components of the detergent, such as surfactants, builders, enzymes, fragrances, dyes, fabric brighteners, fabric softener components, and anti-redeposition polymers and fabric care polymers can occur, and are thought to result in some cases in a reduction in the stain removal and bleaching performance of sodium hypochlorite, or reduction in the aesthetics or detergency performance of the detergent formulation used. Such complex interactions are commonly observed when the stain removal and whitening performance of sodium hypochlorite solutions is measured in the presence of different detergent formulations, whereby some detergent formulations are more compatible with the sodium hypochlorite solution than others. For example, fabric brighteners are often less effective or ineffective when contacted with the bleaching agent of the bleach solution.\nIn the case of manual washing of laundry, where sodium hypochlorite is employed, the relative volume of water to fabrics may be reduced, and hence the levels of suspended soils, both particulate and oily, may be significantly higher than in an automatic washing machine.\nSodium hypochlorite solutions designed to be diluted in such laundering processes may contain a variety of adjuvants to enhance the stability of the hypochlorite or enhance the aesthetics of use of the product (e.g., lower odor). Fragrances, hydrotropic materials (e.g., such as sodium xylene sulfonates), buffers, certain electrolytes (e.g., alkali metal halides, etc.) may be added to hypohalite solutions for such purposes.\nHowever, variation in the overall consumer-perceivable fabric whitening and stain removal still exists when hypohalite solutions are used in combination with commercial detergent formulations. Thus, there continues to be a need for liquid compositions including hypohalite species, which compositions would exhibit improved or more consistent cleaning, bleaching, and whitening performance when used with various modern detergent formulations, or when used under conditions of high loads of suspended soils."}
{"text": "1. Field of the Invention\nThe present invention relates to a nonvolatile semiconductor memory device which allows reusing and a method of reusing same.\n2. Description of the Related Art\nThere is a known nonvolatile semiconductor memory device that has two charge accumulation portions in a memory cell, and has a memory capacity of two bits per memory cell, by storing logic data (“0” or “1”) in each of the two charge accumulation portions (see Japanese Patent Kokai No. H09-91971 (Patent Literature 1) and Japanese Patent Kokai No. 2010-27162 (Patent Literature 2)). The memory cell used in this type of semiconductor memory device has an nMOSFET structure for example, and two charge accumulation portions spaced apart from each other are formed on the drain terminal side and the source terminal side respectively. Storing of data, one bit per charge accumulation portion, two bits per memory cell, can be performed by assigning a state that electric charge is accumulated in the charge accumulation portion to data “0” for example, and assigning a state that electric charge is not accumulated in the charge accumulation portion to data “1”, for example.\nWriting of data into or reading of data from the memory cell of this type and erasing of data are performed in the following method, for example.\nWhen performing the writing of data “0” into the charge accumulation portion on the drain terminal side for example, a positive voltage is applied to the drain terminal and the gate terminal, and the source terminal is held at the ground voltage. Consequently, hot electrons are injected into the charge accumulating portion on the drain side and held therein, thereby writing of data “0” is performed.\nNext, when reading the data stored in the charge accumulation portion on the drain terminal side, a positive voltage is applied to the source terminal and the gate terminal, and the drain terminal is held at the ground voltage. In this instance, a relatively large read current is obtained if electric charge is not accumulated in the charge accumulation portion on the drain terminal side, that is, data “1” is recorded in the charge accumulation portion on the drain terminal side. Conversely, if electric charge not accumulated in the charge accumulation portion on the drain terminal side, that is, data “0” is recorded in the charge accumulation portion on the drain terminal side, the read current becomes smaller relative to the current obtained when data “1” is recorded, due to the effect of the accumulated electric charge. Thus, a difference in magnitude of the read current appears depending on the presence/absence of electric charge in the charge accumulation portion, so that the reading of data can be performed by judging the largeness/smallness of the read current.\nWhen erasing the data stored in the charge accumulation portion on the drain terminal side, a positive voltage is applied to the drain terminal, zero or a negative voltage is applied to the gate terminal, and the source terminal is held in an open state. Consequently, hot holes produced near the drain area are injected into the charge accumulation portion, and the data is erased by the neutralization of the electric charge accumulated in the charge accumulation portion.\nThere exists a use of writing a program to this type of semiconductor memory device and mounting the semiconductor memory device in an appliance. In some cases, the program may be repeatedly rewritten by a user as the need arises. In such event, security issues can arise as discussed in Japanese Patent Kokai No. 2009-146008 (Patent Literature 3).\nTo cope with the security issue, a shipping mode is conceivable in which data, e.g., a program, is written into an EEPROM (Electrically Erasable and Programmable Read Only Memory) of which the capability for writing is then limited. However, if such a memory device in which data is already written has remained unused by a user, the memory device will become just useless for the user.\nConventionally, in order to cope with a situation that such “useless” memory devices having data written beforehand remain unused and are maintained by a user, a business plan has been considered, that includes procedures of buying the unused memory devices from users, erasing data written on them, writing necessary data subsequently, and shipping them to users."}
{"text": "The present invention relates to a ventilated item of clothing such as a coat, jacket, overcoat or the like.\nIt is known that man wears clothing and shoes to protect his body from atmospheric agents such as snow, rain, wind and in particular from the cold.\nThe human body is protected mainly by resorting to various layers of clothing; the first one, underwear, is in direct contact with the body and is in turn covered by successive layers according to the outside temperature and to the environmental conditions.\nAccordingly, it is sufficient to add or remove one or more layers of clothing to achieve comfort and find oneself at an optimum temperature.\nThe human body is inherently provided with xe2x80x9cmechanismsxe2x80x9d which help it to adapt thermally to the environment in which it is placed.\nIn the presence of overheating, for example, the body reacts by increasing perspiration which, by evaporating, allows a natural reduction in body temperature.\nThe heat generated by the human body, in addition to producing perspiration, is also dissipated externally by radiation.\nThis heat, which is always present, warms the air contained between the body and the item of clothing; the air, by rising, causes further overheating and discomfort, for example at the shoulders, which constitute regions of accumulation.\nIf water vapor is unable to escape from the protective enclosure that surrounds the human body (clothing), humidity increases until the vapor condenses and returns to the liquid state of perspiration, soaking the clothing starting from the underwear that constitutes the first inner layer.\nThis unpleasant drawback can currently be remedied by removing the wet item of clothing to replace it with a dry one, for example immediately after the end of a challenging mountain climb, but this causes a sudden cooling of the body and the danger of pneumonia and other chill-related diseases.\nWhile on the one hand protection of the human body against the most adverse cold conditions is very effective by way of the use of highly insulating materials, on the other hand one cannot avoid noting the inability to allow the body to perspire naturally, ensuring the escape of the water vapor produced by perspiration.\nEvidently, during the warm season the problem is more significant and forces many people to take several showers and change clothing continuously through the day.\nAn attempt has been made to solve these drawbacks by using clothing which has special vapor-permeability characteristics, for example by resorting to a material known commercially by the trademark xe2x80x9cGore-Texxe2x80x9d owned by the company W. L. Gore Ass. Inc.; however, such clothing is able to expel only a fraction, often a small fraction, of the vapor produced by sweating and generated by the human body, especially at the regions that are richer in sweat glands, and in any case is unable to ensure an effective changing of the air inside the item of clothing.\nVapor permeation in fact occurs to a limited extent, because a partial pressure of vapor sufficient to expel the sweat (in the vapor phase) outward does not form inside the layer of the item being worn.\nIn other cases, a remedy has been attempted by forming in the items of clothing openings which can be closed to varying extents at the regions where sweat concentrates most, for example under the armpits, but even this does not ensure particular effects, since actual air changing is not produced.\nIt should also be noted that the attempt to increase effectiveness by providing a larger number of openings also has not yielded satisfactory results.\nIn practice, in fact, some portions of the items of clothing always cling directly to the body, particularly the shoulders and chest, and therefore the water vapor generated by the evaporation of body sweat remains trapped between the body and the regions of the clothing that do not cling directly to the body (generally the abdominal region, the lumbar region of the back and most of all the region under the armpits), thus preventing its escape.\nThe aim of the present invention is to provide an item of clothing which solves the drawbacks noted above of conventional ones.\nWithin this aim, an important object of the present invention is to provide an item of clothing which ensures adequate air changing inside the protective enclosure that surrounds the body.\nAnother object of the invention is to allow, for all practical purposes, the natural thermoregulation of the human body.\nThis aim and these and other objects which will become better apparent hereinafter are achieved by an item of clothing, such as a coat, jacket, overcoat or the like, comprising:\nspacer means which are arranged at least in an internal part of the shoulders of the item of clothing so as to generate an interspace between the shoulders of the user and the fabric which the item of clothing is made of;\nelements for venting warm air, produced by the body of the user, that has risen in said interspace."}
{"text": "One example of a hearing device is a hearing-aid device. Modern hearing-aid devices, when employing different hearing programs (typically two to four hearing programs, also termed audiophonic programs), permit their adaptation to varying acoustic environments or scenes. The idea is to optimize the effectiveness of the hearing-aid device for the hearing-aid device user in all situations.\nThe hearing program can be selected either via a remote control or by means of a selector switch on the hearing-aid device itself. For many users, however, having to switch program settings is a nuisance, or it is difficult, or even impossible. It is also not always easy, even for experienced users of hearing-aid devices, to determine, at what point in time which hearing program is suited best and offers optimum speech intelligibility. An automatic recognition of the acoustic scene and a corresponding automatic switching of the program setting in the hearing-aid device is therefore desirable.\nThe switch from one hearing program to another can also be considered a change in a transfer function of the hearing device, wherein the transfer function describes signal processing within the hearing system. The transfer function may depend on one or more parameters, also referred to as transfer function parameters, and may then be adjusted by assigning values to said parameters.\nThere exist several different approaches to the automatic classification of acoustic surroundings. Typically, the methods concerned involve the extraction of different characteristics from an input signal. Based on the so-derived characteristics, a pattern-recognition unit employing a particular algorithm makes a determination as to the attribution of the analyzed signal to a specific acoustic environment.\nAs examples for classification methods and their application in hearing systems, the following publications shall be named: WO 01/20965 A2, WO 01/22790 A2 and WO 02/32208 A2.\nFurthermore, EP 1 670 285 A2, published on Jun. 14, 2006, shall be mentioned, which discloses a training mode for classifiers in hearing devices. It is disclosed that in said training mode, a sound source can be separated by narrow beam-forming. This will isolate the targeted source and, as far as said training mode is on, the classifier will be trained for the targeted source, while other sources of sound are suppressed by said narrow beam-forming. The training provides the classifier with considerable amounts of data on the class represented by the targeted source. This way, an improved reliability of the classification can be achieved.\nNot in all situations, hearing program change based on the classification result provides for an optimum hearing sensation for the user. It would be desirable to provide for an improved basis for choosing a hearing program to switch to and/or for the point in time when to switch hearing programs."}
{"text": "1. Technical Field\nThe disclosure relates to a driving device, and more particularly to a light emitting device (LED) driving device.\n2. Related Art\nGenerally, LED driving device usually first coverts an AC voltage to a DC voltage, and then drives LEDs and controls the brightness of LEDs using a stable DC voltage or a stable DC current.\nHowever, to turn on two or more than two LEDs connected in series, LED driving device must supply the voltage greater than the totality of the turned-on voltages of two or more than two LEDs. Thus, the power source with higher voltage value is required. This may cause higher power consumption. Moreover, if the current passing through LEDs is unstable, the brightness of LEDs may be affected. Thus, it is desirous to develop a new LED driving device, which may improve LED's illumination quality and have less power consumption."}
{"text": "1. Field of the Invention\nThis invention generally relates to filters, and more particularly, to a filter having multiple piezoelectric thin-film resonators.\n2. Description of the Related Art\nRecently, there has been an increasing demand for filters with compact and light resonators and filters configured by combining these filters as wireless communication equipment that may be typically cellular phones has spread rapidly. In the past, dielectric filters or surface acoustic wave (SAW) filters were mainly used. Nowadays, piezoelectric thin film resonators and filters using these resonators are attracting attention, in which the piezoelectric thin film resonators are characterized in that they gave good high-frequency performance and are compact and producible as a monolithic device.\nThe piezoelectric thin-film resonators may be categorized into an FBAR (Film Bulk Acoustic Resonator) type and an SMR (Solidly Mounted Resonator) type. Japanese Patent Application Publication No. 60-189307 and Japanese Patent Application Publication No. 2004-200843 disclose piezoelectric thin-film resonators of FBAR type. FIG. 1 show a conventional FBAR type resonator and a conventional SMR type resonator in parts (a) and (b), respectively. A lower electrode 13, a piezoelectric film 14 and an upper electrode 15 are laminated in this order on a substrate 11 having a cavity 16 or an acoustic multilayer film. The cavity 16 or an alternative acoustic multilayer film is formed below an overlapping portion in which the upper electrode 15 and the lower electrode 13 overlap each other across the piezoelectric film 14. Such a portion is referred to as a resonance portion 23. The cavity 16 in the FBAR may be formed by dry or wet etching from the backside of the substrate 11, which may be made of silicon. The cavity 16 may be formed at a position between the lower electrode 13 and the substrate 11 by removing a sacrificed layer provided on the surface of the silicon substrate 11 by wet etching. In this case, the cavity 16 may be called a gap. The acoustic multilayer film in the SMR has first layers having a relatively low acoustic impedance and second layers having a relatively high acoustic impedance so as to be laminated alternately to a thickness equal to λ/4 where λ is the wavelength of an acoustic wave of the resonator.\nThe upper and lower electrodes may be made of aluminum (Al), copper (Cu), molybdenum (Mo), tungsten (W), tantalum (Ta), platinum (Pt), ruthenium (Ru), rhodium (Rh), iridium (Ir) or the like. The piezoelectric thin film may be made of aluminum nitride (AlN), zinc oxide (ZnO), lead zirconium titanate (PZT), lead titanate (PbTiO3) or the like. The substrate 11 may be made of glass other than silicon.\nThe piezoelectric thin-film resonator of FBAR or SMR type has a loss caused by leak acoustic waves 30 propagated outwards from the resonance portion 23, as shown in parts (a) and (b) of FIG. 1. Portions that are located further out than the resonance portion 23 are referred to as non-resonance portions. The leak acoustic waves 30 propagated through the non-resonance portions are not converted into electric signals and cause a loss. This phenomenon is called lateral leakage of the acoustic waves 30. The filters with the piezoelectric thin-film resonators are required to have restrained lateral leakage and reduced loss. Further, the filters of this type are required to have high resistance to electrostatic discharge (ESD) damage."}
{"text": "The present invention relates generally to a method and system for providing calling party identification and more particularly to a method and system for providing calling party identification for an unidentified calling party to a subscriber connected to the Internet via their telephone.\nMany telephone customers subscribe to caller ID service. Caller ID provides to the subscriber an identification of the calling party. Some calling parties' identity may be unavailable, based upon the location of the calling party or if the calling party selectively blocks the calling party identification. For this case, some telephone customers also subscribe to a feature in which incoming calls which do not include calling party identification (either unavailable or blocked) are first “screened.” For screening, the telephone system first prompts the calling party to audibly identify itself to the telephone system. This audible identification is then played to the subscriber, who can then decide to accept or reject the call.\nCaller ID can also be provided to a telephone subscriber logged onto the Internet via the subscriber's line: If a new call is made to the subscriber's line while connected to the Internet, the subscriber sees a “pop-up” screen on their computer display informing them of the new call and presenting several options for disposal of the call. These dispositions include: taking the call over the Internet, sending the call to the subscriber's voicemail, rejecting the call, and hanging up the Internet connection and accepting the call over the phone. This feature works well when the identity of the calling party is known; however, when the calling party identity is unknown (e.g., the calling party information is unavailable or blocked), the information presented to the subscriber is simply “unknown call.” The subscriber does not have enough information to determine whether to accept or reject the call. Even if the subscriber subscribes to both the caller ID and audible identification feature, the pop-up screen for a call with no calling party information would simply show that the calling party identification is not available. The audible identification cannot be provided to the subscriber."}
{"text": "The invention relates to a central locking installation for a motor vehicle, having a controller and a remote actuating device, which installation is provided with at least one transmitter, operating as a key, and one receiver, operating as a lock, upon actuation of the transmitter for the purpose of locking or unlocking the door locks, a code word and a transmitter recognition signal are transmitted, and whereupon a control pulse is transmitted by the receiver to the controller which causes the locking or unlocking process. The receiver contains a coding device, by means of which a transmitter allocation is made according to the transmitter recognition signal.\nSuch a central locking installation is disclosed in German Patent Specification 3,244,049. In this case, a quantity of defined, different code words (stock of code words) are stored in the transmitter and in the receiver, which for their part are arranged in a sequence. When a control element is actuated, the transmitter transmits one of these code words, to be precise switching on through the sequence, starting with the first code word. A decoding device in the receiver correspondingly converts a received code word into an unlocking signal and supplies this to the controller. If the stock of code words is exhausted, a start is made at the first code word again, cyclically. If an embodiment having a plurality of transmitters is involved, but having different quantities of code words from transmitter to transmitter, the decoding device for processing these code words is equipped with a plurality of decoding channels corresponding to the number of transmitters. When the transmitter is actuated, a fixedly set recognition signal (transmitter recognition signal), depending on the transmitter, is transmitted which differs from the recognition signals of other transmitters and switches on the corresponding decoding channel in the receiver.\nIn the known central locking installation, both the transmitter recognition signal and the number of decoding channels, and hence the number of transmitters used in the central locking installation, must be known in advance. These data must have been set at the manufacturer's works, for example by programming. The user is thus confined to a predetermined number of keys from the start."}
{"text": "Research and development in the field of metals, slags and the like conventionally involve the use of bulky, expensive and relatively unsophisticated melting equipment.\nIn laboratories, melting furnaces are used which are small only as compared to full scale industrial furnaces. They must be permanently installed and normally comprise a steel shell lined with substantially the same refractory used by the industrial furnaces.\nUsing such a refractory lining and with furnace temperatures the same or possibly higher than those of many industrial furnaces, the furnace linings of laboratory furnaces must be made with the lining thicknesses of the industrial furnace linings, to maintain the outside of the furnace at tolerable temperatures for the operating technicians.\nA small portable furnace unit which can be used to melt only a few hundred grams of material under precisely controlled conditions, with a finger-touchable outside surface temperature, is an attractive concept in connection with improving the technology of experimenting and developing high temperature melting materials. The development of such a small furnace unit has been inhibited by the prior art concepts concerning melting furnace constructions."}
{"text": "The present invention relates generally to the manufacture of high speed, high performance MOS semiconductor devices fabricated on strained lattice semiconductor substrates, and MOS devices obtained thereby. Specifically, the present invention relates to an improved method of manufacturing MOS devices including gate insulator layers comprised of high-k dielectric materials on strained lattice semiconductor substrates, which method substantially eliminates, or at least minimizes, stress relaxation of the strained lattice semiconductor layer attendant upon gate insulator layer formation.\nRecently, there has been much interest in investigating the feasibility of various approaches having the aim or goal of developing new semiconductor materials which provide increased speeds of electron and hole flow therethrough, thereby permitting fabrication of semiconductor devices, such as integrated circuit (IC) devices, with higher operating speeds, enhanced performance characteristics, and lower power consumption. One such material which shows promise in attaining the goal of higher device operating speeds is termed xe2x80x9cstrained siliconxe2x80x9d.\nAccording to this approach, a very thin, tensilely strained, crystalline silicon (Si) layer is grown on a relaxed, graded composition Sixe2x80x94Ge buffer layer several microns thick, which Sixe2x80x94Ge buffer layer is, in turn, formed on a suitable crystalline substrate, e.g., a Si wafer or a silicon-on-insulator (SOI) wafer. Strained Si technology is based upon the tendency of the Si atoms, when deposited on the Sixe2x80x94Ge buffer layer, to align with the greater lattice constant (spacing) of the Si and Ge atoms therein (relative to pure Si). As a consequence of the Si atoms being deposited on the Sixe2x80x94Ge substrate comprised of atoms which are spaced further apart than in pure Si, they xe2x80x9cstretchxe2x80x9d to align with the underlying lattice of Si and Ge atoms, thereby xe2x80x9cstretchingxe2x80x9d or tensilely straining the deposited Si layer. Electrons and holes in such strained Si layers have greater mobility than in conventional, relaxed Si layers with smaller inter-atom spacings, i.e., there is less resistance to electron and/or hole flow. For example, electron flow in strained Si may be up to about 70% faster compared to electron flow in conventional Si. Transistors and IC devices formed with such strained Si layers can exhibit operating speeds up to about 35% faster than those of equivalent devices formed with conventional Si, without necessity for reduction in transistor size.\nAnother tactic for improving the performance of semiconductor devices, e.g., MOS devices such as PMOS and NMOS transistors and CMOS devices, involves increasing the capacitance between the gate electrode and the underlying channel region within the semiconductor substrate. Typically, the capacitance is increased by decreasing the thickness of the gate dielectric layer, typically an oxide layer such as a silicon oxide, to below about 100 xc3x85. Currently, silicon oxide, e.g., SiO2, gate dielectric layer thicknesses are approaching about 40 xc3x85 or less. However, the utility of silicon oxide as a gate dielectric is severely limited at such reduced thicknesses, e.g., due to direct tunneling through the gate dielectric layer to the underlying channel region, thereby increasing the gate-to-channel leakage current and an increase in power consumption.\nInasmuch as further reduction in the silicon oxide gate dielectric thickness is impractical in view of the above-mentioned increase in gate-to-channel leakage current, various approaches have been investigated for reducing the gate-to-channel leakage current while maintaining a thin SiO2 xe2x80x9cequivalent thicknessxe2x80x9d, i.e., the thickness of a non-SiO2 dielectric layer determined by multiplying a given SiO2 thickness by the ratio of the dielectric constant of the non-SiO2 dielectric to that of SiO2, i.e., knon-SiO2/kSiO2. Thus, one approach which has been investigated is the use of materials with dielectric constants higher than that of silicon oxide materials as gate dielectric materials, whereby the xe2x80x9chigh-kxe2x80x9d dielectric materials, i.e., materials with dielectric constants of about 5 or above, replace the conventional silicon oxide-based xe2x80x9clow-kxe2x80x9d dielectric materials with dielectric constants of about 4 or below. The increased capacitance k (or permittivity ∈) of the gate dielectric material advantageously results in an increase in the gate-to-channel capacitance, which in turn results in improved device performance. Since the capacitance C is proportional to the permittivity ∈ of the gate dielectric material divided by the thickness t of the gate dielectric layer, it is evident that the use of a high-k (or high-∈) material permits use of thicker gate dielectric layers, i.e.,  greater than 40 xc3x85, whereby both greater capacitance and device speed are obtained with less gate-to-channel leakage current.\nTypically, high-k dielectric materials, i.e., with kxe2x89xa75, suitable for use as gate dielectric layers in the manufacture of semiconductor devices, are formed with a physical thickness from about 40 to about 500 xc3x85, typically 40-100 xc3x85 (or a SiO2 equivalent thickness less than about 40 xc3x85), and comprise metal and oxygen-containing material including at least one dielectric material selected from the group consisting of metal oxides, metal silicates, metal aluminates, metal titanates, metal zirconates, ferroelectric materials, binary metal oxides, and ternary metal oxides. Suitable metal oxides include aluminum oxide, hafnium oxide, zirconium oxide, lanthanum oxide, titanium oxide, tantalum oxide, tungsten oxide, cerium oxide, and yttrium oxide; suitable metal silicates include zirconium silicate, and hafnium silicate; suitable metal aluminates include hafnium aluminate and lanthanum aluminate; suitable metal titanates include lead titanate, barium titanate, strontium titanate, and barium strontium titanate; suitable metal zirconates include lead zirconate; and suitable ferroelectric and/or ternary metal oxides include PST (PbScxTa1xe2x88x92xO3), PZN (PbZnxNb1xe2x88x92xO3), PZT (PbZrxTi1xe2x88x92xO3), and PMN (PbMgxNb1xe2x88x92xO3). Deposition of the high-k metal oxide layers and/or post-deposition treatment of the high-k metal oxide layers typically involves processing at elevated temperatures, e.g., at about 500-900xc2x0 C. in the case of aluminum oxide (Al2O3) deposition from an AlCl3/O2 ambient.\nHowever, an important concern in the manufacture of practical semiconductor devices utilizing strained semiconductor layers, e.g., strained Si layers, is the requirement for maintaining the tensilely strained condition of the strained semiconductor layer throughout device processing, without incurring significant strain relaxation disadvantageously leading to reduction in electron/hole mobility resulting in degradation in device performance characteristics. For example, many device fabrication steps, including for example, the above-described high-k dielectric deposition and post-treatment, frequently involve high temperature processing at temperatures on the order of about 900-1,100xc2x0 C. for intervals sufficient to result in significant relaxation of the tensile strain of the Si layer, which in turn, results in a lowering of the electron and hole mobilities therein to values comparable to those of conventional Si layers, whereby the potential advantages attributable to enhanced electron/hole mobility in the strained Si layer are partially or wholly lost.\nAccordingly, there exists a need for improved methodology for fabrication of semiconductor devices with strained semiconductor layers, notably strained Si layers, which substantially eliminates, or at least minimizes, deleterious stress relaxation during device processing at elevated temperatures, e.g., as in the formation of high-k dielectric gate insulator layers as part of a process sequence for the manufacture of MOS-type transistors and CMOS devices.\nThe present invention, wherein processing for deposition of high-k dielectric gate insulator layers and buffer layers therefor, e.g., metal oxide-based high-k dielectrics and silicon oxide-based buffer layers, forming part of a sequence of steps for fabricating MOS-type transistors and CMOS devices, is performed at temperatures which effectively eliminate, or at least minimize, disadvantageous strain relaxation of the strained lattice semiconductor arising from the thermal annealing. As a consequence, the inventive methodology facilitates manufacture of high speed, high performance, reduced power consumption semiconductor devices utilizing strained semiconductor technology. Further, the methodology afforded by the present invention enjoys diverse utility in the manufacture of numerous and various semiconductor devices and/or components therefor which require use of strained semiconductor technology for enhancement of device speed and lower power consumption.\nAn advantage of the present invention is an improved method for manufacturing a semiconductor device comprising a strained lattice semiconductor layer.\nAnother advantage of the present invention is an improved method for manufacturing a semiconductor device comprising formation of a high-k dielectric gate oxide layer on a strained lattice semiconductor layer without incurring significant stress relaxation.\nStill another advantage of the present invention is an improved method for fabricating a MOS-type semiconductor device on a strained lattice semiconductor layer.\nYet another advantage of the present invention are improved MOS-type semiconductor devices fabricated on strained lattice semiconductor substrates.\nAdditional advantages and other aspects and features of the present invention will be set forth in the description which follows and in part will become apparent to those having ordinary skill in the art upon examination of the following or may be learned from the practice of the present invention. The advantages of the present invention may be realized and obtained as particularly pointed out in the appended claims.\nAccording to the invention, the foregoing and other advantages are obtained in part by a method of manufacturing a semiconductor device, comprising the sequential steps of:\n(a) providing a semiconductor substrate comprising a strained; lattice semiconductor layer at an upper surface thereof, the strained lattice semiconductor layer having a pre-selected amount of lattice strain therein;\n(b) forming a thin buffer/interfacial layer of a low-k dielectric material on the upper surface of the semiconductor substrate; and\n(c) forming a layer of a high-k dielectric material on the thin buffer/interfacial layer of a low-k dielectric material, wherein:\nsteps (b) and (c) are each performed at a minimum temperature sufficient to effect formation of the respective dielectric layer without incurring, or at least minimizing, strain relaxation of the strained lattice semiconductor layer.\nAccording to embodiments of the present invention, step (a) comprises providing a semiconductor substrate including a crystalline, graded composition Sixe2x80x94Ge layer, with a lattice-matched crystalline silicon (Si) layer on a first side of the Sixe2x80x94Ge layer and comprising the strained lattice semiconductor layer; and according to particular embodiments of the present invention, step (a) further comprises providing a semiconductor substrate including a crystalline Si layer on a second, opposite side of the Sixe2x80x94Ge layer;\nIn accordance with embodiments of the present invention, step (b) comprises forming the thin buffer/interfacial layer of a low-k dielectric material having a dielectric constant k less than 5 and at a thickness from about 2 to about 6 xc3x85.\nAccording to particular embodiments of the present invention, step (b) comprises forming the thin buffer/interfacial layer of a low-k dielectric material from at least one material selected from the group consisting of silicon oxides and silicon oxynitrides and at a temperature ranging from about 200 to about 400xc2x0 C., by means of an atomic layer deposition (ALD) method selected from chemical vapor deposition (CVD), molecular beam deposition (MBD), and physical vapor deposition (PVD).\nEmbodiments of the invention include performing step (c) by forming a layer of a high-k dielectric material having a dielectric constant k greater than 5 and a thickness from about 40 to about 100 xc3x85.\nAccording to particular embodiments of the present invention, step (c) comprises forming the layer of a high-k dielectric material from at least one metal and oxygen-containing material selected from the group consisting of metal oxides, metal silicates, metal aluminates, metal titanates, metal zirconates, ferroelectric materials, binary metal oxides, and ternary metal oxides, i.e., step (c) comprises forming the layer of a high-k dielectric material from at least one material selected from the group consisting of aluminum oxide, hafnium oxide, zirconium oxide, lanthanum oxide, titanium oxide, tantalum oxide, tungsten oxide, cerium oxide, yttrium oxide, zirconium silicate, hafnium silicate, hafnium aluminate, lanthanum aluminate, lead titanate, barium titanate, strontium titanate, barium strontium titanate, lead zirconate; ferroelectric oxides, ternary metal oxides, PST (PbScxTa1xe2x88x92xO3), PZN (PbZnxNb1xe2x88x92xO3), PZT (PbZrxTi1xe2x88x92xO3), and PMN (PbMgxNb1xe2x88x92xO3); wherein step (c) comprises forming the layer of a high-k dielectric material at a temperature ranging from about 200 to about 400xc2x0 C., by means of an atomic layer deposition (ALD) method selected from chemical vapor deposition (CVD), molecular beam deposition (MBD), and physical vapor deposition (PVD).\nIn accordance with embodiments of the present invention, the method further comprises sequential steps of:\n(d) forming an electrically conductive layer on the layer of a high-k dielectric material; and\n(e) patterning the electrically conductive layer, the layer of a high-k dielectric material, and the thin buffer layer of a low-k dielectric material to form at least one gate insulator layer/gate electrode stack on at least one portion of the upper surface of the semiconductor substrate;\n(f) implanting dopant species of one conductivity type into the semiconductor substrate utilizing the at least one gate insulator layer/gate electrode stack as an implantation mask, thereby forming at least one pair of shallow depth source/drain extension regions in the semiconductor substrate vertically aligned with opposite side edges of the at least one gate insulator layer/gate electrode stack;\n(g) forming insulative sidewall spacers on the opposite side edges of the at least one gate insulator layer/gate electrode stack;\n(h) implanting dopant species of the one conductivity type into the semiconductor substrate utilizing the at least one gate insulator layer/gate electrode stack with the insulative sidewall spacers thereon as an implantation mask, thereby forming at least one pair of deeper source/drain regions in the semiconductor substrate vertically aligned with opposite side edges of the sidewall spacers; and\n(i) thermally annealing the thus-formed structure for a minimum interval sufficient to activate the dopant species implanted in the at least one pair of shallow depth source/drain regions and in the at least one pair of deeper source/drain regions without incurring, or at least minimizing, strain relaxation of the strained lattice semiconductor layer.\nAccording to particular embodiments of the present invention, step (i) comprises performing laser thermal annealing (LTA) or rapid thermal annealing (RTA) at a temperature from about 1,200 to about 1,400xc2x0 C. for from about 1 to about 100 nanosec.\nAnother aspect of the present invention are PMOS or NMOS transistors and CMOS and IC devices manufactured according to the above method.\nAdditional advantages and aspects of the present invention will become readily apparent to those skilled in the art from the following detailed description, wherein embodiments of the present invention are shown and described, simply by way of illustration of the best mode contemplated for practicing the present invention. As will be described, the present invention is capable of other and different embodiments, and its several details are susceptible of modification in various obvious respects. Accordingly, the drawings and description are to be regarded as illustrative in nature, and not as limitative."}
{"text": "Employee benefit plans are widely used today by both employees and employers to provide for an individual employee's retirement needs. Pursuant to these employee benefit plans, employees can establish financial accounts under the benefit plan. These financial accounts and the employee benefit plans are established and structured to comply with various tax and employment laws.\nThese benefit plans are administered by a plan administrator who works with the employer (the plan sponsor) to help establish and maintain the benefit plan. The plan administrator is a financial institution or a third party administrator that operates to administer the employee benefit plan. Among other responsibilities, the plan administrator will maintain detailed records showing all transactions in the individual accounts, and will, pursuant to the individual account holder instructions, cause transactions to be executed in the individual's financial account. In some cases the plan administrator may be part of a larger financial institution that includes a broker dealer, and in other cases the plan administrator will be a separate entity which will work with a broker dealer. A typical scenario would include a broker dealer establishing a trust account which holds all of the assets for a retirement plan, but the broker dealer typically does not maintain records for each of the individual accounts. The plan administrator would maintain records as for each of the individual accounts and would convey instructions to the broker dealer for specific transactions which are to be executed in the trust account. The plan administrator would maintain records of all the transactions for each of the individual accounts. Collectively, the combination of all the individual accounts would correspond to the holdings in the trust account held by the broker dealer.\nTypically, benefit plans will be established to provide individual employees with a range of different investment options representing a broad spectrum of different possible asset classes. Frequently, these investment options include a variety of different mutual funds, and other fixed income investments, and money market funds. Company stock may also be an option for some plans.\nIndividual employees participating in these retirement plans can elect to have a certain percentage of their paycheck or a flat dollar amount deposited into their financial account established under the benefit plan. Frequently, the employer will also make contributions to the employee's financial account under the benefit plan. In connection with opening an account under the benefit plan the plan administrator will receive detailed personal information for the individual employee. This information can include the individual's age, date of hire, number of dependents, annual salary, job title/profession, home mailing address, social security number, e-mail address, and possibly a personal identification number (PIN). The plan administrator maintains a data storage system tracking all of the investments for each of the financial accounts under the benefit plan. The plan administrator also causes different transactions to be executed, such as the buying and selling of securities, pursuant to instructions from the individual employees.\nIn many instances the individuals will make their own decisions regarding the buying and selling of different securities. The plan administrator operates a computer system which can track all of the transactions and deposits, withdrawals and earnings for the individual financial accounts. Further, the individuals can generally choose to either increase or decrease the amount of their paycheck which they direct to their financial account.\nOn a periodic basis, which is generally quarterly, the plan administrator will issue an account statement to the individual employees who have individual financial accounts established pursuant to a benefit plan. These account statements typically include information about the financial account such as the value of the account at the beginning of the present calendar year, the value of the account at the time statement was prepared, the assets held in the account, and a description of all transactions in the account during the account period.\nFIG. 1 shows an overview of the architecture of a system 100 for implementing account statements for a benefit plan under the prior art. This system includes a data storage system 110 for storing financial account information. The data storage system also includes information identifying the individual owner, or individual associated, with a financial account, and includes personal information for the individual. On a periodic basis typically quarterly a processor system 112 extracts data from the data storage system 110 and then uses this data to prepare an account statement. In some cases the processor system may include multiple processor systems and the plan administrator may outsource the preparation and delivery of the account statements. In such situation, the plan sponsor would provide account information to a third party and the third party would use a processor system to format and prepare the information so that the account statement could be prepared and delivered to the individual associated with the account. This account statement could be delivered to the individual by printing the account statement and sending it via mail 114, or the account statement could be formatted for transmission via e-mail 116, or the account statement could be formatted so that it could be made accessible to the individual over the web 118 using a personal computer with an internet browser. Numerous methods have been developed for using computer systems such as the Internet to securely deliver confidential information. For example, the following US Patents and International patent applications, which are incorporated herein in their entirety by reference, provide discussion of details of delivering confidential information using the internet: U.S. Pat. Nos. 6,367,010 B1 (METHOD FOR GENERATING SECURE SYMMETRIC ENCRYPTION AND DECRYPTION); 6,014,688 (E-MAIL PROGRAM CAPABLE OF TRANSMITTING, OPENING AND PRESENTING A CONTAINER HAVING DIGITAL CONTENT USING EMBEDDED EXECUTABLE SOFTWARE); 6,304,897 B1 (METHOD OF PROCESSING AN E-MAIL MESSAGE THAT INCLUDES A REPRESENTATION OF AN ENVELOPE); and international patent applications: WO 98/49643; WO 01/03367; WO 99/34286; WO 00/46729; WO 01/78491.\nThe value of the account statement to the individual is that it gives the individual some idea of the past performance of the account, and lets the individual review recent transactions in the account. In the past these account statements have not provided individualized financial planning advice regarding steps that an individual could take to improve the performance of the account relative the individuals retirement needs.\nIn the past when an individual wanted to receive specific advice regarding financial planning for a retirement account, the individual could go to a number of different sources. Traditionally, an individual would enlist the services of a financial adviser. The financial advisor would review the individual's retirement needs and make recommendations for changes to be made in connection with a financial account.\nIn recent years a number of software tools have been made available to individuals, so that they can determine there own financial plan without going directly to a financial advisor. Two example's of software tools which are available include a retirement planning tool provided by mPower Advisors LLC, San Francisco, Calif. (website www.mpower.com) and a retirement planning tool provided by Financial Engines, Palo Alto, Calif. (website www.financialengines.com). Financial Engines is shown as an assignee on a number of patents related to retirement planning tools including U.S. Pat. Nos. 6,292,787 (“Enhancing utility and diversifying model risk in a portfolio optimization framework”); 6,125,355 (“Pricing module for financial advisory system”); 6,021,397 (“Financial advisory system”); 6,012,044 (“User interface for a financial advisory system”); 5,918,217 (“User interface for a financial advisory system”), which are incorporated herein by reference.\nAnother system has been proposed where an individual can receive financial investment advice such that, they are directed to investments in specific trusts, or investment vehicles, which are managed pursuant to certain asset allocation models, as discussed in U.S. Pat. No. 6,154,732 “System for providing investment advice and management of pension assets” which is incorporated herein by reference.\nIn general operation the prior retirement planning tools appear to require that an individual seek out advice. This means, for example, that an individual most go to either Mpower or Financial Engines website, or other website that provides access to a financial planning tool, and input information into the financial planning tool. This information could include an individuals age, current income, financial assets held in a retirement account, age at which an individual would like to retire etc. Based on this information, the financial planning tool can then provide an estimate showing the likely amount of money that one would have upon retirement. Further, the financial planning tool can also make recommendations as to changes in investments, and future financial planning steps to improve the performance of the financial account to improve the probability that an individual will have sufficient assets at the time of retirement.\nOne disadvantage with the present system is that the individual must take affirmative steps to get financial planning advice. Specifically, an individual must make time available to seek out financial planning advice. This means an individual must locate financial planning tools or resources and then use the financial planning tool or resource to get financial planning advice. Frequently, individuals do not make the time to seek out financial planning advice.\nOne disadvantage with the present account statements is that they typically report historical information, and they do not provide advice as to future financial planning strategy. Basically, these account statements provide a snapshot of where the account stands at the time of the account statement, and some information regarding the past condition of the account."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for fabricating a dynamic random access memory (DRAM) cells and, more particularly, to a DRAM cells provided with a word line on bit line (WOB) type memory cell including a buried bit line structure.\n2. Discussion of Related Art\nA metal-oxide-semiconductor (MOS) type DRAM has a memory cell comprising one MOS transistor and one capacitor connected thereto. As developments have been made in DRAM device technique to achieve high integration and high speed response, each capacitor is shrunken to be of such a reduced size that the amount of charges stored in the capacitor is decreased. The decrease in the amount of charges results in soft errors that may destroy the content of memories. To overcome this problem, a method for increasing the occupied area of each capacitor has been proposed, in which storage nodes composed of polysilicon are formed in a semiconductor substrate in order to increase capacitance of the capacitor. The stacked capacitor according to this method is disposed on a transfer gate transistor and connected to a source or drain of the transfer gate transistor. Bit lines of the DRAM cell are normally composed of metal lines and disposed on an interlayer insulation layer over the word lines. The bit lines are connected to the source and drain regions of the transfer gate transistor through contact holes in the interlayer insulation layer (or a passivating insulation layer).\nFIG. 1 is a plan view of a memory cell array of a DRAM according to a related art.\nReferring to FIG. 1, in the surface of a semiconductor substrate 1 are formed a plurality of word lines 17a, 17b, 17c and 17d which run parallel with one another in rows, a plurality of bit lines 55 which run parallel with one another in columns, and a plurality of memory cells (MC) arranged at the respective intersections of the word lines and the bit lines. Each memory cell comprises one transfer gate transistor 53 and one capacitor 64. The transfer gate transistor 53 comprises a pair of source/drain regions 46 and 46 formed in the surface of the semiconductor substrate 1, and gate electrodes (word lines) 17b and 17c formed between the source/drain regions 46 and 46 with a gate insulation layer 15 interposed therebetween. A thick insulation layer is formed on the gate electrodes 17b and 17c. Subsequently, contact holes 29, 28 and 29 are formed in a predetermined portion of the insulation layer so as to reach the source/drain regions 46 and 46 of the transfer gate transistor 53.\nReference numerals 29 and 29 denote capacitor node contact portions and reference numeral 28 denotes a bit line contact portion. The contact holes 29, 28 and 29 formed by photolithography and the etching method are gap-filled with the plug of a conductive layer such as a doped polysilicon layer. The plug of the conductive layer is brought in contact with the semiconductor substrate 1 in the source and drain regions 46 and 46. A bit line contact hole 91 located over an element isolating insulation layer. In the memory cell array, the word lines have a predetermined width and arranged in parallel with a predetermined spacing from one another. The plural bit lines 55 running in parallel with one another in columns are isolated with an interlayer insulation layer or the like on the plural word lines 17a, 17b, 17c and 17d.\nNow, manufacturing steps of the DRAM memory cell shown in FIG. 1 will be described with reference to sectional views of FIGS. 2A to 2L.\nFIGS. 2A to 2L are cross-sectional views taken along the line A-A' of FIG. 1 showing the manufacturing steps of the DRAM memory cell.\nReferring to FIG. 2A, an element isolating insulation layer 11 and a channel stopper region (not shown) are formed in predetermined regions on the main surface of a P-type semiconductor substrate 1. A gate insulation layer 15, a polysilicon layer 17 and an interlayer insulation layer 19a are sequentially formed on the surface of the semiconductor substrate 1.\nThe element isolating insulation layer 11 may be formed by a selective oxidation method such as a LOCOS (Local Oxidation of Silicon) method or other methods including STI (Shallow Trench Isolation). The gate insulation layer 15 is formed by the thermal oxidation method. The polysilicon layer 17 and the interlayer insulation layer 19a are each deposited to a thickness of 1000-2000 .ANG. by the CVD method.\nReferring to FIG. 2B, word lines 17a, 17b, 17c and 17d are formed by photolithography and the etching method. The interlayer insulation layer 19a of the patterned oxide film is left on the surface of the word lines 17a-17d.\nReferring to FIG. 2C, an insulation layer is formed on the whole surface of the semiconductor substrate 1 by the CVD method, and is etched by an anisotropic reactive ion etching (RIE) to form a sidewall spacer 20 on the peripheries of the word lines 17a-17d. Impurity ions 40, arsenic are implanted under an implantation energy 30 KeV, a dose of 4.0.times.10.sup.15 /cm.sup.2 in the surface of the silicon substrate 1 by using the word lines 17a-17d covered with the insulation layer 19a and the spacer 20 as masks to form the source and drain regions 46 and 46 of the transfer gate transistor.\nReferring to FIG. 2D, the surface of the semiconductor substrate 1 is planarized with an interlayer insulation layer 26a e.g., a BPSG (Borophosphorosilicate Glass) film. Contact holes 31 and 33 are formed in the bit line contact portion 28 and the capacitor node contact portion 29 by photolithography and the etching method, which is followed by deposition of doped polysilicon. Then, polysilicon plugs 28 and 29 are formed in the contact holes by an etch-back method.\nThe plugs may be formed not only by the etch-back technique using the RIE but also by other methods including CMP (Chemical Mechanical Polishing).\nReferring to FIG. 2E, an insulation layer 61 is deposited on the whole surface of the semiconductor substrate 1, isolating the plugs 28 and 29. Contact holes (not shown) are formed over the bit line contact portion 28. A conductive layer such as a doped polysilicon layer or a metal layer, and a metal suicide layer, etc. are formed on the surface of the semiconductor substrate 1, which are patterned by photolithography and the etching method. As a result, the bit lines (not shown) are formed.\nSubsequently, an etching stopping layer 63 such as a nitride (Si.sub.3 N.sub.4) film having a thickness of more than 100 .ANG. is formed and then a silicon oxide (SiO.sub.2) film 65a having a thickness of more than 5000 .ANG. is formed on the surface of the nitride film 63.\nHere, the bit lines are disposed above the element isolating insulation layer of the memory cell array and, in the perpendicular direction to the word lines, arranged in parallel with the active region of the memory cell array having two transfer gate transistors as MOS transistors formed thereon.\nReferring to FIG. 2F, a capacitor isolating layer 65 for isolating the adjacent capacitors is formed by patterning the oxide film 65a by the etching method. The selective ratio of the etching of the nitride film 63 which is an etching stopping layer to the oxide film 65a is extremely high. Therefore, in this etching step, the nitride film 63 is etched at a rate different from that of the oxide film 65a.\nReferring to FIG. 2G, contact holes 70 and 70 are formed so as to reach the plug of the capacitor node contact portion 29 on the source and drain regions 46 and 46 by photolithography and the etching method.\nReferring to FIG. 2H, a polysilicon layer 72 of a thickness of 500-1500 .ANG. is deposited on an inner surface of the contact hole 70, on the surface of the nitride film 63 and on the surface of the capacitor isolating layer 65 by the CVD method. Then, a thick resist 75 is applied over a surface of the polysilicon layer 72.\nReferring to FIG. 2I, the resist 75 is etched back to expose a part of the polysilicon layer 72.\nReferring to FIG. 2J, the exposed surface of the polysilicon layer 72 is selectively removed by anisotropic etch or the like. As a result, the polysilicon layer 72 is isolated on the surface of the capacitor isolating layer 65 to form a lower electrode 80 of the capacitor.\nReferring to FIG. 2K, the resist 75 is removed by the etching method and furthermore the capacitor isolating layer 65 is removed by a plasma etching method. Then, a dielectric layer 84 such as a nitride film is formed on the surface of the lower electrode 80.\nReferring to FIG. 2L, an upper electrode 85 at a thickness of 2000-3000 .ANG. of a polysilicon layer is formed on the surface of the dielectric layer 84 by the CVD method. Thereafter, an insulation layer 88 and an interconnection layer 90 are formed to complete the manufacturing steps of the memory cell of the DRAM.\nThe memory cell of the DRAM according to the related art has the arrangement in which the plugs used as the capacitor node contact portion and the bit line contact portion are formed in the source and drain regions of the transfer gate transistor by photolithography and the etching method after the word lines are formed. These patterns in a linear arrangement are aligned to dispose between the transfer gates. This requires that the transfer gates be spaced farther apart than is otherwise desired to provide ample room between transfer gates for making desired contacts to active areas. Such increase in transfer gate pitch works against the integration of the memory cell. Another problem lies in that soft errors occur due to a parasitic capacitance, since the capacitor of memory cell adjoins the peripheries and the top portion of the bit lines."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nIn an embedded system, there are typically constraints on hardware resources thereby limiting processing capabilities. This limits the speed in which the system is able to load all the necessary software components to become useable. In particular, embedded systems with a graphical user interface (GUI) will take an inordinate amount of time as GUIs are typically resource intensive.\nTherefore, an unaddressed need exists in the art to address the aforementioned deficiencies and inadequacies."}
{"text": "This invention relates to the sorting, counting and grouping of conveyed items, particularly wrapped packages containing one or more articles, such as cookies or biscuits.\nFlat articles such as biscuits, cookies or the like are conventionally packed individually or in small groups in hose-like wrappers. These packages have projecting, welded flaps at two oppositely located edges. Thereafter, the individual, elemental packages are, in exactly counted quantities, united into larger, combined packages. Since not all the packages delivered by a wrapping machine are in an acceptable condition (occasionally empty wrapping sleeves are discharged by the wrapping machine or two or more successive packages are not entirely severed from one another at their flaps), they have to be carefully monitored and the unacceptable ones separated before they are counted and gathered into combined packages.\nHeretofore it has not been possible to perform the monitoring, separating, counting and grouping of the elemental packages in a fully automatic manner. These operations had to be performed in part manually by labor which has caused a substantial increase in the manufacturing costs."}
{"text": "Various elements and structures have been used with electronic devices to provide electrical paths between connection elements on the devices. The present invention is directed to novel uses and applications of and improved processes for making air bridge structures configured to convey electrical signals on an electronic device."}
{"text": "The present invention relates to computer readable storage media for repeatedly contacting software and/or software end-users and providing information, such as software passwords, marketing, advertising, and/or promotional material.\nSoftware developers are often victims of illicit copying and unauthorized use of their software in violation of contractual obligations imposed by licensing agreements and subject to civil and criminal penalties under various domestic and foreign laws. Unauthorized entities range from a relatively small percentage of the total users to an overwhelming majority of illegal users. Such unauthorized use not only amounts to theft of the developers\"\" intellectual property, but also reduces the number of programs sold and therefore the associated profitability of the developer. This may ultimately diminish the creative effort expended by the software developers due to the reduced financial incentive. The advent of the Internet has contributed to the proliferation of pirated software, known as xe2x80x9cwarezxe2x80x9d, which is easily located and readily downloaded.\nVarious strategies have been employed to make unauthorized duplication and use of software more difficult. One such approach is to provide a hardware xe2x80x9ckeyxe2x80x9d which is typically installed in the parallel port of the computer to provide a software interlock. If the key is not in place, the software will not execute. This method is relatively expensive for the developer and cumbersome for the authorized user while remaining vulnerable to theft by duplication of the hardware key.\nAnother approach requires the user to enter a serial number or customer identification number during installation of the software. Missing or invalid registration information prevents installation of the software. This approach is easily defeated by transferring the serial number or customer identification number to one or more unauthorized users.\nYet another approach requires registering the software with the manufacturer or distributor to obtain an operational code or password necessary for installation of the software. Again, once the operational code or password is obtained, it may be perpetually transferred along with pirated copies to numerous unauthorized users.\nVarious copy protection strategies have been employed to reduce the number of unauthorized copies available. This approach is generally disfavored by users who may have a legitimate need to make backup or archival copies or transfer a copy to a new computer or hard drive.\nWhile prior art strategies have enjoyed various levels of success in reducing unauthorized use of software, they often impose a significant burden on the authorized users or are easily defeated by unauthorized users. As such, software developers need an apparatus and/or method for reducing unauthorized use of software which does not burden the authorized users to dissuade them from purchasing and using the protected software.\nThus, one object of the present invention is to provide an apparatus and method for improving software security throughout the lifetime of the software.\nAnother object of the present invention is to provide an apparatus and method for monitoring the number of users of a software product, both authorized and unauthorized.\nYet another object of the present invention is to provide an apparatus and method for the software manufacturer to maintain contact with the user over an extended period of time to provide repeated opportunities for promotional and marketing purposes, for example.\nA further object of the present invention is to provide a method and apparatus for reducing unauthorized use of software which facilitate periodic software updates and forwarding of information, when and if desired.\nA still further object of the present invention is to identify those entities responsible for unauthorized copying or use of software so that appropriate action may be taken, such as disabling the software, requesting payment from the user, or seeking civil or criminal penalties.\nAnother object of the present invention is to provide a method and apparatus for reducing unauthorized software use which deactivates unauthorized copies when an unauthorized user attempts to obtain a password.\nIn carrying out the above objects and other objects, features, and advantages of the present invention, a method for providing continuing contact with a software user to repeatedly transfer information to the user includes receiving registration information from the user, transferring authorization information to authorize continued use of the software beyond a current authorization interval, and transferring marketing, promotional, or advertising information each time the user or software requests authorization for continued use beyond the current authorization interval. In one embodiment, the present invention transfers authorization information which authorizes use of the software for an authorization interval based on value of the software with value being determined by retail price or the availability of subsequent versions, for example. The user or software may automatically contact a manufacturer or authorized agent preferably prior to expiration of the current authorization interval to provide continuous uninterrupted operation of the software.\nThe present invention contemplates, but does not require, more frequent password updates for more complex software because it is generally more costly to develop (and therefore more valuable to users) whereas less costly software would require fewer password updates to reduce administrative costs associated with password maintenance.\nPassword or authorization code updates may be obtained automatically or manually. Automatic updates are accomplished using electronic communication between the manufacturer\"\"s computer (or an authorized representative) and the user\"\"s computer. Updates may be performed by a direct modem connection, via email, a web browser, or the like. The particular time and nature of updates and the user interface utilized to implement the updates may vary by manufacturer or product. Manual updates are performed by advance or periodic notifications generated by the software to alert the user that password updates are required or will be required in the near future. The user may then contact the manufacturer for the specific password update via telephone, mail, email, or the like. Password advisories normally occur prior to the periodic termination of the operating period which may be measured by program starts, elapsed running time, calendar period, etc. Password updates may be in the form of alphanumeric and/or encrypted passwords or of any other conventional type.\nPreferably, the user must provide registration information prior to receiving the original or updated password or authorization code. Registration information may be entered by the user or automatically acquired (and transmitted for automatic updates) by the software. Registration information may include a serial number, registration number, TCP/IP address, user name, telephone number, computer specific information, etc. This information may be encoded and/or encrypted to make it less susceptible to tampering by unauthorized users. The registration information is preferably monitored and compared to previously captured information to control the number of authorized copies of software and/or identify unauthorized users. If unauthorized use is suspected, a password or authorization code may be provided which subsequently disables the software, either immediately or after some period of time so that an authorized user is provided an opportunity to rectify the information which caused deactivation.\nA number of advantages result from various implementations of the present invention. For example, the present invention reduces unauthorized use of software without imposing a significant burden on authorized users. The present invention controls the number of copies of authorized software by monitoring registration information and deactivation of suspected pirated copies. Requiring authorized users to periodically update a password or authorization code provided by a password administrator improves accuracy of contact information for marketing related products and distribution of product updates. The present invention also provides a variable level of software security which can be tailored to the particular application depending upon the value of the application to potential software pirates. By providing an apparatus and method for the software manufacturer to maintain contact with the user over an extended period of time, the present invention provides repeated opportunities for promotional and marketing purposes, for example. This feature of the invention facilitates periodic forwarding of information, such as advertising, marketing, and promotional information, when and if desired.\nThe above advantages and other advantages, objects, and features of the present invention will be readily apparent from the following detailed description of the best mode for carrying out the invention when taken in connection with the accompanying drawings."}
{"text": "The present invention relates to slow release tablets for oral administration and, more particularly, to tablets which permit a prolonged period of contact between the medication and the buccal and gingival mucosa as well as the hard and soft tissues of the mouth.\nMedications are generally administered with a pharmaceutical carrier with specific desired characteristics which can, for example, influence the rate of release of the medication from the carrier. The choice of carrier is thus determined by the desired medical outcome. For example, medications which have local effects in the mouth and throat are often administered in both solid and liquid form, depending upon the desired pharmaceutical effect.\nLiquid and semi-liquid dosage forms, such as mouthwashes and toothpastes are frequently used for oral hygiene, but have the disadvantage of relatively short contact periods with the hard and soft tissues of the mouth and throat. Solid dosage forms, such as troches and lozenges, offer certain advantages for such local administration of oral medications. The terms lozenge and troche are used synonymously for any form of tablet designed to be held in the mouth for slow dissolution and release of medication such that prolonged contact of the medication with the mouth and throat is ensured. The medications administered with troches or lozenges are primarily local anesthetics, antiseptics, astringents or anti-tussives. Lozenges may be made by fusion, a candy molding process or by compression. Troches are usually manufactured by compression, as are most other tablets.\nAlthough currently available troches and lozenges have the advantage of enabling medication to be in prolonged contact with the mouth and throat, such tablets still have a number of disadvantages. First, although troches and lozenges do not dissolve immediately, the rate of release of medication is still relatively rapid, on the order of about 15 minutes for total release. Second, these dosage forms have not been specifically tested with herbal medications, or medications derived from botanical materials. Thus, the efficacy of these dosage forms with herbal medications is unknown. Hereinafter, the term \"herbal medication\" refers to a medication derived from botanical materials or a biologically active extract of these materials.\nThere is thus a widely recognized need for, and it would be highly advantageous to have, a solid dosage form for oral administration of medication with local effects on the mouth and throat, which permits release of the medication over a prolonged period of time, which has been specifically tested for use with herbal medications and which permits prolonged contact between the active ingredient of the medication and the hard and soft tissues of the mouth and throat."}
{"text": "The state of the art is well summarized in a report prepared by Henry Persson and entitled \"Design, Equipment and Choice of Tactics are Critical When Fighting Large Tank and Bund Fires.\" That report is further identified as Brandforsk Project: 612-902; Swedish National Testing and Research Institute, Fire Engineering, SP Report 1992:02. The purpose of that report was to develop knowhow based on the experience and recommendations of fire experts to provide a practical basis for the design and planning of foam extinguishing systems for large tank and bund fires. A bund fire is one within the embankment or dike surrounding a storage tank.\nA traditional approach to the fighting of such fires has been to direct streams of water and foam onto the fire site through monitors or even hand-held nozzles. In order to successfully extinguish large fires using traditional techniques it is necessary to have available an adequate supply of water and foam concentrate to allow the application of foam liquid at a minimum rate of 6.5 l/m.sup.2 /min to the burning surface for some 60 to 90 minutes. The report indicates agreement among the experts that concentrating foam application on as small an area in the tank as possible is far superior to the previously accepted technique of fighting tank fires with several small monitors distributed around the circumference of the tank. Concentrating the foam application at one point more quickly establishes a bridge head, or initial foam cover, thus increasing the effectiveness of subsequently applied foam.\nAmong other findings of the report are that no successful fire extinguishing operation has been verified in tanks over about 45 m in diameter. Indeed, some experts hold that a tank of 45 m in diameter represents about the largest that can be extinguished with mobile equipment. It appears to be the general consensus of the experts that tanks up to at least 60 m in diameter can be extinguished if the tanks are equipped with fixed \"over-top\" systems to apply foam. It is considered possible as well to extinguish fires in even larger tanks if the over-top system is supplemented with a bottom feed system.\nA fixed, over-top system comprises permanently installed piping and foam sprinkler nozzles within the tank itself at a level above the liquid surface when the tank is filled to capacity. A bottom feed system employs a hose array with foam deploying nozzles adapted to float on the surface of the stored liquid and to rise and fall with the liquid as the tank is filled and emptied. Both systems require connection to a water source and to a supply of foam concentrate. That connection may be permanent one through a direct attachment to the water mains and to a store of foam but the systems are more commonly supplied from a mobile unit which is connected to a system through hoses at the time of need. Both systems are difficult to maintain and are essentially impossible to test without contamination of the tank contents.\nFires aboard tankers carrying either crude oil or refined petroleum products pose many of the same problems as do fires in stationary tanks. Consequently, tanker fires have been traditionally fought using much the same tactics used in the fighting of stationary tank fires. However, the difficulties of access and of the coordination of equipment, personnel and decision-making are ordinarily vastly greater in a tanker fire than those encountered at land locations.\nFires in tankers and stationary storage tanks, while relatively uncommon, pose enormous risks. Those risks include the threat of injury or death to people aboard the ship or in the area or engaged in fighting the fire, the likelihood of huge property losses, and the nearly certain contamination of soils, beaches, ground and surface water and air. Further, the intense thermal radiation always threatens to ignite adjacent structures and tanks thus compounding the risks and increasing the potential losses.\nWith this background, it can readily be appreciated that fire fighting tactics and systems which can more quickly and surely bring under control and extinguish fires aboard tankers and in tanks, particularly those fires in large stationary or mobile tanks, is of great environmental and economic importance."}
{"text": "1. Field of the Invention\nThis invention relates to a confocal laser scanning microscope.\n2. Related Background Art\nA scanning optical microscope for deflecting a light beam to scan a specimen, which has acquired high resolving power but still retains convenience is known (Japanese Patent Laid-Open Publication No. 219919/1986).\nThis scanning optical microscope has the drawback that stray reflected and scattered light produced when the laser beam passes through the relatively large number of optical elements lowers the measuring precision of the microscope."}
{"text": "This invention relates to novel polymer product compositions, polymerization feed compositions and monomer purification therefor. In particular, it relates to a crosslinked, high modulus, high impact strength, thermoset polymer of polycycloolefin units which is formed via a metathesis-catalyst system.\nDistillation is commonly used in the preparation of polycycloolefins, for example dicyclopentadiene. Various purities of dicyclopentadiene are available. The invention preferably uses economically available purified dicyclopentadiene which is readily polymerized. Purified dicyclopentadiene forms a substantially cross-linked thermoset polymer when polymerized as disclosed herein.\nAny good thermoset polymer should meet at least two criteria. It should have desirable physical properties and it should lend itself to easy synthesis and forming. Among the most desirable physical properties for many polymers is a combination of high impact strength and high modulus. A standard test for impact strength is the notched Izod impact test, ASTM No. D-256. For an unreinforced thermoset polymer to have good impact strength, its notched Izod impact should be at least 1.5 ft. lb./in. notch. It is desirable that this good impact strength be combined with a modulus of at least about 150,000 psi at ambient temperature. Thermoset polymers with high impact strength and high modulus find useful applications as engineering plastics in such articles of manufacture as automobiles, appliances and sports equipment. Among the critical factors in the synthesis and forming of a thermoset polymer are the conditions required to make the polymer set up or gel. Many thermoset polymers require considerable time, elevated temperature and pressure, or additional steps after the reactants are mixed before the setting is complete.\nA thermoset homopolymer having high impact strength and high modulus has been described by Klosiewicz in U.S. Pat. Nos. 4,400,340, 4,469,809 and 4,436,858 (with plasticizer) and by Leach in 4,458,037 (a foam). Characteristics of thermoset polymers include insolubility in common solvents such as gasoline, naphtha, chlorinated hydrocarbons, and aromatics as well as resistance to flow at elevated temperatures.\nWork has been done on the metathesis copolymerization of dicyclopentadiene with one or more other monomers to produce soluble copolymers. This copolymer formation has resulted in the production of unwanted insoluble by-products. U.S. Pat. No. 4,002,815, for instance, which teaches the copolymerization of cyclopentene with dicyclopentadiene, describes an insoluble by-product and suggests that the by-product could be a gel of a dicyclopentadiene homopolymer.\nSome other work, usually in an attempt to produce soluble polymers, has been done on the metathesis polymerization of dicyclopentadiene. Japanese unexamined published patent applications KOKAI Nos. 53-92000 and 53-111399 disclose soluble polymers of dicyclopentadiene. Several syntheses of soluble polymers of dicyclopentadiene have produced insoluble by-products. Takata et al, J. Chem. Soc. Japan Inc. Chem. Sect., 69, 711 (1966), discloses the production of an insoluble polymerized dicyclopentadiene by-product from the Ziegler-Natta catalyzed polymerization of dicyclopentadiene; Oshika et al, Bulletin of the Chemical Society of Japan, discloses the production of an insoluble polymer when dicyclopentadiene is polymerized with WCl.sub.6, AlEt.sub.3 /TiCl.sub.4 or AlEt.sub.3 /MoCO.sub.5 ; and Dall Asta et al, Die Makromolecular Chemie 130, 153 (1969), discloses an insoluble by-product produced when a WCl.sub.6 /AlEt.sub.2 Cl catalyst system is used to form polymerized dicyclopentadiene.\nIn U.S. Pat. No. 3,627,739 ('739), dicyclopentadiene is gelled with unactivated catalyst and then heated for an hour.\nPampus et al in U.S. Pat. No. 3,873,644 discloses a method of producing graft polymers from a cyclic olefin to obtain thermoplastic products of high impact strength.\nOfstead in U.S. Pat. Nos. 4,020,254 and 3,935,179 disclose metathesis polymerization of cycloolefins to obtain a rubbery polymer. Streck et al in U.S. Pat. Nos. 3,974,092 and 3,974,094 discloses a catalyst for preparation of polyalkenamers. Each discloses polyalkenamers of low reduced melt viscosity. Percent gel obtained using the system of Streck et al is very low. A typical percent gel is about 2 to 6% with 14 being the highest value disclosed. Wilkes in U.S. Pat. No. 3,084,147 discloses thermalpolymerization of dicyclopentadiene. The thermal polymerization is carried out in a nonpolymerizable solvent at about 500.degree. to 550.degree. F.\nStafford in U.S. Pat. No. 3,446,785 discloses a polymerization of olefins. The polymerization discloses the production of a product which is a viscous liquid polymer or a brittle, semi-solid polymer.\nNutzel et al in U.S. Pat. No. 3,684,787 discloses preparation of polyalkenamers. The polymers can be isolated by pouring the solution of polymer into 3 to 5 times its quantity of a solution of a lower alcohol in which an age resister is dissolved. Alternatively, the solution can be introduced into boiling water and the solvent removed with steam. Mensel states in column 1, line 19 that crosslinked products are of no industrial interest. The polymers obtained have a rubber-like character.\nVergne et al in U.S. Pat. No. 3,557,072 discloses plastomers derived from dimethanooctahydronaphthalene and their method of manufacture. Amorphous polymers are obtained in the form of a hard white mass which is washed with methanol and then ground and dried to a white powder.\nTenney et al and Tenney alone in U.S. Pat. Nos. 4,136,247, 4,136,248, 4,136,249 and 4,178,424 disclose ring-opened cycloolefin polymers and copolymers. These polymer products are thermoplastics which can be thermoformed.\nMinchak and Minchak et al respectively in U.S. Pat. Nos. 4,002,815 and 4,380,617 each disclose polymerization of cycloolefins to form polymers which may be isolated by precipitation using an alcohol or by steam or hot water stripping. The polymers produced have inherent viscosities from about 0.1 to about 10 and are greater than 90% soluble in solvent. Viscous cements and plastic polymers solidify in from about 30 minutes to 180 minutes. The reaction is short stopped in less than 2 hours by addition of an alcohol. Minchak in U.S. Pat. No. 4,426,502 discloses bulk polymerization of cycloolefins by reaction injection molding in less than about 2 minutes using an organoammonium molybdate or tungstate catalyst.\nDeWitt et al in U.S. Pat. No. 4,418,179 discloses impact modification cycloolefins by polymerization using an organoammonium molybdate or tungstate catalysts in the less than 2 minutes using reaction injection molding.\nOshika et al in the Bulletin of the Chemical Society of Japan, line 41, pages 211-217 (1968) discloses ring opening polymerization of norbornene and its derivatives by MoCl.sub.5, WCl.sub.6 and ReCl.sub.5 catalysts. Dark crude polymer is obtained which is dissolved and reprecipitated with methanol.\nU.S. Pat. No. 4,002,815 discloses the use of a metathesiscatalyst system which employs a dialkylaluminum iodide, an alkylaluminum diiodide or a mixture of trialkylaluminum compounds with elemental iodine to produce substantially gel-free copolymers of cyclopentene and dicyclopentadiene.\nU.S. Pat. No. 4,069,376 discloses the use of a three component catalyst comprised of a soluble tungsten compound, a dialkylaluminum chloride or alkylaluminum dichloride, and a dialkylaluminum iodide or alkylaluminum diiodide to produce substantially gel-free norbornene-dicyclopentadiene copolymers.\nU.S. application Ser. No. 526,835 filed Aug. 26, 1983 and assigned to the same assignee, discloses a cellular crosslinked poly(dicyclopentadiene) which is made with a metathesis-catalyst system. The cellular polymer is made by injecting the catalyst system, which includes an alkylaluminum activator, into a reaction vessel which is preheated, preferably to a temperature from about 100.degree. C. to about 125.degree. C.\nNot only is it desirable that the thermoset polymer have high impact strength, but it is also desirable that it be easily synthesized and formed. A reaction injection molding (sometimes hereinafter referred to as RIM) process achieves this second goal by in-mold polymerization. The process involves the mixing of two or more low viscosity reactive streams. The combined streams are then injected into a mold where they quickly set up into a solid infusible mass. For a RIM system to be of use with a particular polymer, certain requirements must be met: (1) the individual streams must be stable and must have a reasonable shelf-life under ambient conditions; (2) it must be possible to mix the streams thoroughly without their setting up in the mixing head; (3) when injected into the mold, the materials must set up to a solid system rapidly; and (4) any additives-fillers, stabilizers, pigments, etc. - must be added before the material sets up. Therefore, the additives selected must not interfere with the polymerization reaction.\nInterpenetrating polymer networks (IPNs) are a type of polymer alloy consisting of two (or more) crosslinked polymers. They are more-or-less intimate mixtures of two or more distinct, crosslinked polymer networks held together by permanent entanglements with few, if any covalent bonds between the polymers. The entanglements in IPNs must be of a permanent nature and are made so by self-crosslinking of the two polymers. They are introduced either by swelling a crosslinked polymer with monomer and the crosslinking agent of another polymer, and curing the swollen polymer in situ or by mixing the linear polymer, prepolymers, or monomers in some liquid form solution, or bulk, together with crosslinking agents, evaporating the vehicle (if any), and curing the component polymers simultaneously.\nIPNs possess several interesting characteristics in comparison to normal polyblends. Formation of IPNs is the only way to intimately combine crosslinked polymers, the resulting mixture exhibiting (at worst) only limited phase separation. Normal blending or mixing of polymers results in a multiphase morphology due to the well-known thermodynamic incompatibility of polymers. However, if mixing is accomplished simultaneously with crosslinking, phase separation may be kinetically controlled by permanent interlocking of entangled chains.\nMolecular and supramolecular perspectives are believed to be important in considering interpenetrating networks. On linked polymer networks intermeshed through molecular chain segment entanglements. These are totally non-separable without chain breaking (i.e. degradation). A semi-interpenetrating network (SIPN) then is a crosslinked network having a non-crosslinked polymer dispersed through it on a molecular scale. These in theory are separable without bond breaking. Neither the interpenetrating network nor the SIPN represents a heterophase polymer structure in a morphological or supramolecular sense.\nThe supramolecular scale microscopically observable heterophases are present SIPN and IPN. In this context an interpenetrating network is any system of two polymers where the supramolecular structure involves two intermeshed continuous polymer phases irrespective of whether or not they are crosslinked.\nBecause molecular mixing of dissimilar uncrosslinked polymers is thermodynamically unfavorable (for entropic reasons) it is unusual to form a stable, molecular IPN from two uncrosslinked polymers. They generally will separate into a heterophase polymer system which might or might not be a supramolecular IPN.\nThe forming of supramolecular interpenetrating networks of two uncrosslinked polymers is called polymer blends. \"Co-reacted\" polymers is a process which results in block copolymers and sometimes results in heterophase polymer systems which could exhibit supramolecular IPN morphology, depending on block length."}
{"text": "Amyotrophic lateral sclerosis (ALS) is a fatal neurodegenerative disorder caused by motor neuron death (Rowland et al, N. Engl. J. Med., 2001, 334, 1688-1700) and characterized in part by the presence of abnormal aggregates of insoluble protein in selectively vulnerable populations of neurons and glia. ALS, an orphan disease, is estimated to afflict about 87,000 people worldwide, but its prevalence would be much higher were it not for the fact that ALS patients survive for only 3 to 5 years on average after diagnosis. Approximately 10% of ALS cases are familial, with the rest of the cases being sporadic (Rowland et al., N. Engl. J. Med., 2001, 334, 1688-1700). Approximately 20% of the cases of familial ALS are caused by inherited mutations in the protein Cu/Zn superoxide dismutase (SOD1) (Rosen et al, Nature, 1993, 362, 59-62). Rodent models in mutant SOD are often used as a disease model because of its phenotypic and pathologic resemblance to sporadic and familial human ALS (Dal Canto et al., Brain Res., 1995, 676, 25-40; Wong, et al., Neuron, 1995, 14, 1105-1116; Bruijin et al., Science, 1998, 281, 1851-1854; Bruijn et al., Neuron, 1997, 18, 327-338; Wang et al., Hum. Mol. Genet., 2003, 12, 2753-2764; Wang et al., Neurobiol. Dis., 2002, 10, 128-138; Jonsson, et al., Brain, 2004, 127, 73-88).\nThe causes of sporadic ALS remain unknown, and the clinical courses are variable, suggesting that multiple factors are involved. Different hypotheses have been proposed, such as glutamate-mediated excitotoxicity, impaired mitochondrial function, oxidative stress, and aberrant protein aggregation (Dib et al., Drugs, 2003, 63, 289-310; Strong et al., Pharmacology & Therapeutics, 2003, 98, 379-414; Kunst et al., Am. J. Hum. Genet., 2004, 75, 933-947; Bruijn, et al., Annu. Rev. Neurosci, 2004, 27, 723-749; Dibernardo et al., Biochimica et Biophysica Acta, 2006, 1762, 1139-1149). Riluzole, which decreases glutamate excitotoxicity, is the only FDA approved ALS drug (Jimonet et al., J. Med. Chem., 1999, 42, 2828-2843.). However, it can only extend median survival life for 2 to 3 months, suggesting mechanisms other than glutamate-mediated excitotoxicity should be considered during ALS drug development (Miller et al., ALS and Other Motor Neuron Disorders, 2003, 4, 191-206; Taylor et al., Neurology, 2006, 67, 20-27).\nImbalances in protein homeostasis, which can be caused by cell stress or expression of certain mutant proteins, can result in the appearance of alternative conformational states that are able to self-associate to form protein aggregates and inclusion bodies. In familial as well as sporadic forms of ALS and mutant SOD1 transgenic models, aberrant protein aggregation has been reported as a common feature (Brnijin et al., Science, 1998, 281, 1851-1854; Leigh et al., Cytoskeletalpathology in Motor Neuron Diseases, in: Rowland, L. P. (Editor), Amyotrophic Lateral Sclerosis and Other Motor Neuron Disease, Raven Press, New York, 1991, pp. 3-pp. 23; Mather et al., Neurosci. Lett., 1993, 160, 13-16; Pasinelli et al., Neuron, 2004, 43, 19-30; Watanabe, et al., Neruobiol. Dis., 2001, 8, 933-941). The ubiquity of this feature in various forms of ALS suggests that perhaps drugs that inhibit aberrant protein aggretation may provide new and improved treatment options. Thus, there is an urgent need to identify and develop drugs for the treatment of ALS."}
{"text": "The present invention relates generally to the field of pivoted hand held tools, and more specifically to a ratchet mechanism for adjusting the tension of the blades and for blade replacement of a hand pruner.\nPivoted implements having elongated members disposed for cooperative engagement about a pivotable joint are widely used. In particular tools such as pruning snips generally comprise two elongated members, typically made of stamped or forged metal or other suitable material, disposed for cooperative engagement about the pivotable joint. Typically, each member includes a jaw at the front end portion thereof, an opposed tang, and a pair of handles connected to the tangs.\nAn improvement to the basic hand pruner is disclosed in U.S. Pat. No. 5,697,159 entitled Pivoted Hand Tool to Linden and is incorporated herein by reference. The \"\"159 patent discloses a hand pruner in which one of the handles is rotatable. A transmission mechanism is provided to interconvert the rotation force applied to the rotatable handle into a rectilinear force used to draw the handles together.\nAs a result of wear on the blades and general use of the tool, it may become necessary to replace the blades, or remove them for sharpening and or repair. Additionally, it may be necessary to adjust the tension between the two blades for optimal operation of the pruning tool. However, the removal of the blades or adjustment of the tension requires additional tools that are difficult to carry when working in the field.\nIt would be desirable to provide a pruning tool that could allow for adjustment of the tension of the blades without the need for external tools. Further, it would be desirable to provide a hand pruning tool that could allow the removal of a damaged blade and installation of new blade in the field without the use of additional tools.\nOne embodiment includes a pivoted hand tool including a first and second handle, each handle having a jaw member extending therefrom. A pivot connects the first and second handles. A ratchet is operatively engaged with the pivot and at least one of the handles and selectively rotates the pivot in a clockwise and counter-clockwise direction by movement of the handles toward or away from one another\nIn another embodiment a method for adjusting a pivoted tool includes providing a first and second handle, each handle having a jaw member extending therefrom, the first and second handles being pivotally secured to one another with a threaded pivot threadably engaged with one of the handles. A ratchet assembly is attached to the threaded pivot and the ratchet assembly is linked to the handles. The threaded pivot is tightened by rotating the threaded pivot in one direction relative to the one of the handles by movement of the handles toward or away from one another.\nIn a further embodiment a hand pruner includes a first and second handle. A first and second blade are operatively secured to the first and second handles. A pivot connects the first and second handles. A three position ratchet is operatively engaged with the pivot and coupled to the first handle with a lever. Movement of the first handle toward and away from the second handle pivots a ratchet gear in a first and second direction about the pivot. The ratchet is adjustable between a first position in which movement of the handle toward the second handle rotates the pivot relative to the handles in a first direction, to a second position in which movement of the handle away from the second handle rotates the pivot in a second direction opposite the first direction. The ratchet includes a third neutral position in which movement of the handles toward and away from one another does not rotate the pivot relative to the handles."}
{"text": "1. Field of the Invention\nThis invention relates to the movable seating structures generally called strollers, and more particularly to incorporation of the present invention in a foldable stroller of the type such as disclosed in my U.S. Pat. No. 3,995,882.\n2. Description of the Prior Art\nPrior art folding strollers, such as disclosed in my prior U.S. Pat. No. 3,995,882, have been used for the transport of infants, children, invalids, persons partially paralyzed or suffering from some debilitating disease or condition, wherein such person may be too weak to walk unassisted, and in other desired applications. Often it is desired to have the stroller foldable so as to allow the stroller to be conveniently transported or stored when not in use. Such storage may be, for example, in the trunk of an automobile, under a bed, or the like.\nWhen in use, however, such prior art foldable strollers provide only one position of the seat and back relative to the framework and/or the folding structure. Such single position was not, in certain applications and uses, the most comfortable for the person being transported in the stroller. For example, if the person being so transported in the stroller falls asleep, a more recumbent position may be desired than is provided in the upright seating condition. Additionally, when the stroller is pushed over or down some obstruction, such as a curb or the like, it is desired for both comfort and safety, to shift the center of gravity of the combined person and stroller so that the person does not fall from the stroller or undergo the sensation of being about to fall from the stroller. The seat and back structure must, at all times during use, provide secure seating conditions regardless of the relative movement thereof with respect to the supporting framework. For convenient use, the weight of the entire structure must be kept to a minimum consistent with the safety requirements so that it may be easily carried, moved and placed into and removed from a storage location by one person.\nDefinitions\nIn this specification, the following definition of terms used will be assumed.\nAn \"X member\" shall mean a two-element scissors-like assembly in which the two cross elements are elongated members pivotally attached together at their intermediate regions, the assembly having the general appearance of an \"X\". The elements may lie substantially in a common plane, such as the \"X\" member 14 in FIGS. 1-11 (except where they cross in the middle), or may have their elements on one side of the pivotal attachment lying generally in one plane and their elements on the other side of the pivotal attachment generally lying in a nonparallel plane, such as the \"X\" member 104 in FIGS. 12-16. Due to the pivotal attachment of the two cross elements, when the top or bottom element ends are drawn together horizontally, the \"X\" member lengthens vertically, and vice versa. Likewise, when the left or right element ends are drawn together vertically, the \"X\" member lengthens horizontally, and vice versa.\nThe term \"collapsed\" shall mean structurally reduced in dimension, as by folding and/or the pivoting toward one another of pivotally joined members.\nThe term \"expanded\" shall mean structurally increased in dimension, as by unfolding and/or the pivoting away from one another of pivotally joined members.\nThe terms \"left\", \"right\", \"front\", and \"back\" (or \"rear\") shall refer to those positions as they pertain to the orientation of a person sitting in the mobile seating arrangement.\nThe term \"framework\" shall mean basic mobile seating arrangement structure upon which, or with which, structure implementing features of the present invention can be placed or implemented. Typically, a \"framework\" will include the wheels, or other floor support structure, and associated connection members such as front and rear leg members.\nDescribing elements of construction and functions in a symmetrical apparatus\nDue to the symmetry of the mobile seating arrangement about a front-to-rear center line, the right side counterparts of elements of the construction are not visible in the various left elevation views presented in the accompanying drawing. For convenience in describing such left elevation views, only the elements visible from the left will be described as singular elements. It will be understood that corresponding elements of the apparatus which would only be viewed from a right elevation depiction exist and have corresponding physical and mechanical features and functions. Occasionally, it is more reasonable to refer to symmetrical parts in a plural sense, such as mentioning arm rests, even though the view being described shows only one arm rest. It will be understood, therefore, that where singular terminology is used for elements that obviously have mirror image counterparts, such singular terminology is used for convenience only."}
{"text": "In the architectural industry numerous software tools are currently used throughout the engineering design process, as a means to guarantee accuracy, provide consistency, automate tedious or numerically intensive tasks, and to help visualize the resulting design. In general these tools can be categorized into design, simulation, and visualization tools. Furthermore, geographic information systems (GIS) may be applied to any of the three groups of tools.\nThe simulation tools are generally used to model, forecast, and analyze the movement of vehicles, people, or material goods. In aviation design, the software is used to model customer activities such as passenger movement or baggage handling and studies of landside and airside traffic. Some of the applications in transportation planning include toll plaza evaluation, freeway and corridor studies, environmental impact analysis, light and heavy rail transit studies, and ITS (intelligent transportation systems) assessments. Military applications include analysis of range training facilities and troop movement. In the newly emerging field of homeland security, applications include evacuation planning and emergency management.\nThe visualization tools can be categorized in two areas; 3D modeling and virtual reality (VR) software. The 3D modeling software is used to produce high end, photo realistic quality, images and animated video. The VR software is used to generate immersive environments which allow a user to freely move about and interact within a virtual world.\nBoth types of software are used to visualize engineering designs in numerous application areas. Advances in software and hardware technology now allow the quality of the visual images produced to contain a high degree of photo realism. GIS software provides tools which allows visualization from a geographic perspective. These tools provide the ability to compile data from many sources, to link location and information to that data, and to interpret how it interrelates.\nIn the development of a video or a virtual environment, an artist creates animated objects without reference to any previous data generated by use of simulation software. While a good artist can generate animated objects which may appear to be an accurate representation for the environment and circumstances, the artists do not create an accurate representation using simulation data."}
{"text": "1. Field of the Invention\nEmbodiments relate to a digital image processing apparatus and a method of controlling the apparatus, and more particularly, to a digital image processing apparatus, which displays a setting status necessary for an operation of capturing an image and/or for other operations on a display device to facilitate user's manipulation, and a method of controlling the apparatus.\n2. Description of the Related Art\nGenerally, a digital image processing apparatus includes any apparatus which receives or processes an image or uses an image recognition sensor, such as a film camera, a digital camera, a personal digital assistant (PDA), a phone camera, a personal computer (PC) camera, etc.\nIn particular, the digital image processing apparatus may receive an image through an imaging device, process the image input with a digital signal processor, compress the processed image to generate an image file, and store the image file in a storage medium.\nThe digital image processing apparatus may also display the image input through the imaging device or an image from the image file stored in the storage medium on a display device such as a liquid crystal display (LCD) to allow a user to view the displayed image."}
{"text": "One particular application of the invention relates to use in association with ductwork or piping where it become necessary to service an element in such ductwork or piping. Such servicing generally requires the removal of the element from the ductwork or piping and then subsequent reinstallation of the serviced item or installation of a replacement item. In most cases the most common means for removal or placement of such items comprises the use of a crane and whilst this can be convenient during the construction phase of a plant it is not so convenient once the plant is operating and in fact may not be possible due to the location of the element to be removed and the site.\nA particular application of the invention relates to use with flanged pipes which are interconnected through their flanged connections and/or which accommodate flow control or sensing elements through their flanged connections.\nThe discussion throughout this specification and claims, of the background and prior art to the invention is intended only to facilitate an understanding of the present invention. It should be appreciated that the discussion is not an acknowledgement or admission that any of the material referred to was part of the common general knowledge in Australia and the world as at the priority date of the application."}
{"text": "In train systems, the detection of train movement and direction is used to manage traffic between trains. It is also used to determine location of train units of the train systems in storage locations.\nIn some existing train systems, the detection of train movement is typically performed by a communication-based train control (CBTC) system in which vehicle on-board controllers (VOBC) of each train unit when in operation communicate with the wayside Zone Controller to detect train movement and identify the location of each train in storage locations.\nThere are other existing methods used to detect train movement via a secondary detection system by use of axle counters and track circuits. The secondary detection system is used in train systems that include train units equipped with steel wheels and traveling on steel rails. Track circuits detect when train units shunt the steel rails, and axle counters detect and count steel axles when passing, via the detection head of the secondary detection system.\nSome train systems have non-steel (e.g., rubber) tire systems such as monorails, and the secondary detection system is not suitable for detecting train movement in these train systems. In such cases, some existing methods require simulation of steel wheel-to-steel rail interface so that the secondary detection system may be used. In yet another existing method, a position of the train units in storage locations is identified manually by a physical connection (e.g., a plug) between the train and the storage location. In these existing methods, train movement is detected; however, the direction of the movement is undetectable."}
{"text": "1. Field of the Invention\nThe present invention relates to a data transmission system, a data recording and reproducing apparatus, and a recording medium each having a data structure of error correcting code, a data structure for use in a recording medium, a method for arranging a parity of an error correction code, a method for transmitting information data, a method for receiving information data, a method for recording information data in a recording medium, and a method for reproducing information data from a recording medium.\nIn particular, the present invention relates to a data transmission system for transmitting data through a communication line, a data recording and reproducing apparatus for recording data into a recording medium and reproducing data from the recording medium, and a recording medium such as an optical disk, a magneto-optical disk, a magnetic disk, a compact disk, a magnetic tape, a semiconductor memory, or the like, each of the data transmission system, the data recording and reproducing apparatus, and the recording medium having a data structure of error correcting code, a data structure for use in a recording medium, a method for arranging a parity of an error correcting code, a method for transmitting information data, a method for receiving information data, a method for recording information data in a recording medium, and a method for reproducing information data from a recording medium.\n2. Description of the Prior Art\nWhen transmitting and recording a great amount of digital data, an error correcting code (referred to as an ECC hereinafter) is widely used for the purpose of improving the reliability of the data.\nThe ECC word generated by attaching a parity to an original data, and then, the redundancy of the data having the parity increases, however, the ECC has the capabilities of detecting and correcting errors generated in the process of transmitting, recording and reproducing the data.\nIn general, the ECC has a higher error correcting capability when the ECC word has an increased redundancy, i.e., when the ECC word includes a greater amount of parity symbols. Therefore, the ECC word includes a plurality of parity symbols or bytes.\nConventionally, when ECC encoding a great amount of data upon transmitting or recording the data, the data is previously divided into a plurality of blocks each having a predetermined length, and a parity is attached to each block. Then, the parity is arranged at the tail of the block to be corrected, generally. That is, in the conventional process of data transmission, the data has been transmitted with the parity added to the tail of each data block. On the other hand, in the data recording process, the parity has been recorded in a position adjacent to each data block.\nAn arrangement of the data and the parity of the ECC word in the conventional data recording method is concretely shown in FIG. 1.\nFIG. 1 shows a logical data format of one block of an optical disk in which data and a parity are recorded by a prior art method for arranging the parity of the ECC.\nThe block is divided into a data area 90, a row parity area 54, and a column parity area 53, depending on the type of information data recorded. In the present case, the parity is a part of the ECC word. As is apparent from FIG. 1, the row parity area 54 and the column parity area 53 are partially overlapped with each other.\nOne block is comprised of 156 data frames 51 and 12 parity frames. In the present case, the frame 51 means a group of data or parities corresponding to one row of each block, and constructs an ECC word in the row direction.\nEach data frame 51 is comprised of data of 158 bytes and a row parity of 10 symbols each of which consists of one byte (referred to as a C1 parity hereinafter) for error correction of the data. Then, thirteen data frames constitute one sector 52. Therefore, one block is comprised of twelve sectors 52 in this case. At the head of each sector 52 is written a physical sector address 50 of each sector 52 (referred to as a sector address hereinafter). Further, it is noted that a synchronous (Sync) pattern (not shown) is generally added to the head of each frame.\nIn the column parity 53, twelve bytes in the direction of column constitute one parity (referred to as a C2 parity hereinafter), and the parity is used for correcting an error of the 156-byte data series located upward in an identical column. The above-mentioned 156-byte data and 12-byte parity construct an ECC word in the column direction. FIG. 1 shows a logical format of one block, and the block is physically written continuously in a sequence from the frame located in the uppermost position of the block onto a track of an optical disk in a manner as shown in FIG. 2. Therefore, paying attention to the data series in the direction of column, data are discretely arranged on the disk, and this arrangement is referred to as an interleave.\nThus, conventionally, in both of (a) the data arrangement in the direction of row, i.e., the data series which is not interleaved and (b) the data arrangement in the direction of column, i.e., the data series which is interleaved, the data area 90 in which data is arranged and the parity area 91 in which the parity is arranged have been clearly distinguished from each other.\nHowever, according to the recent trend of increasing the storage capacities and reducing the access times of optical disks, upon reading out data from any optical disk in which the data have been written by the prior art arrangement method, there have been such problems as follows.\nAttention is now paid to timings for detecting sector addresses upon continuously reading out information along a track. FIG. 3 shows the timings for the above operation. In this case, the information data is sequentially read out from the frame located in the uppermost position of the block shown in FIG. 1, and therefore, the twelve sector addresses are detected at a predetermined regular time interval. However, since the twelve column parity frames of the column parity 53 are subsequently read out, and therefore, a certain time interval continues for which no sector address is detected continues. Thereafter, the next block is read out, and then, the similar phenomenon occurs repetitively.\nThus, according to the conventional method for arranging the parity of the ECC, the timings for detecting the sector addresses becomes irregular. Therefore, upon retrieving a desired sector address, a complicated control circuit and a complicated calculator circuit are necessary, and this leads to a cause of retarding the time of reading-out. In this case, for example, by providing only the row parity without adding the column parity itself, the timings for detecting the sector addresses can be made regular. However, with only the row parity, the data correcting capability is reduced to impair the reliability of the recorded data."}
{"text": "Various types of conveyors have been utilized for conveying objects in industrial production lines. Objects may be conveyed from work station to work station individually or in groupings, depending on the object and the task to be performed. It may or may not be important to maintain any spacing or control of the objects during some or all of the travel. For example, apples being conveyed may simply be stacked randomly on a conveyor, while bottles being filled may be held rigidly in place by a filling machine that has received the bottles from a conveyor.\nCertain conveyor belts (sometimes also called chains) are made of a plurality of interconnected links, driven by motors that engage the conveyor belt. Such conveying systems are commonly employed in the transportation of manufactured goods and articles, and for containers. With these typical systems, the motor drives a toothed drive sprocket that engages complimenting driving recesses or “dogs” formed on the conveyor belt. These drive units can be disposed in any number along the length of the conveyor belt. Such a drive unit and conveyor system is disclosed in U.S. Pat. No. 6,119,848 which is assigned to the assignee of the present invention, and is incorporated herein by reference in its entirety for all purposes.\nLink type conveyor belts are sometimes designed having a knuckle/socket joint arrangement wherein one part of the link has a rounded knuckle and the opposite part has a socket formed by two extending edges. The knuckle of one link fits into the socket of a neighboring link. The knuckle is able to move in various directions within the socket, which allows for the conveyor system as a whole to curve and move.\nThe interconnected links typically have a platform member connected to or formed integral with the link's upper (conveying) surface. The platform member is generally shaped to match the neighboring platform members on other links such that the links can turn in a plane or twist while moving around curved sections of the conveying system, yet are also shaped such that the cracks and spaces formed between the links are minimized. The platform members can be connected to the links in several different ways. For instance, the platforms may have pegs extending therefrom which match corresponding slots on the links. Alternatively or additionally, the platforms can have snap springs which lock into place on corresponding sections of the links. Such a knuckle link with a platform surface member is disclosed in U.S. Pat. No. 6,209,716 which is owned by the assignee of the present invention and incorporated herein by reference in its entirety for all purposes.\nIn conventional endless loop conveyor systems, the conveyor belt (or chain) follows a path having a conveying section and a return section. In the conveying section the conveyor carries the conveyed objects from a process start point to a process finish point (depending on the application), and in the return section the conveyor returns from the finish point back to the start point. Thus, depending on the orientation of the conveyor, a substantial portion of the conveyor may not be used to actually convey objects at any given point in time. Accordingly, the cost of an endless conveyor per unit of active conveying section is increased due to the manner in which endless conveyors are conventionally configured and driven. The increased costs include original purchase, repair, maintenance, and/or installation costs for the length of the conveyor required along the return path.\nOther conventional systems used to convey objects include “puck” systems, in which a freely moving puck carries a conveyed object along a conveyor. Pucks are often used to hold objects in a given orientation for a certain processing step. If the object is a container, the process may be filling, capping, etc. Guides along the sides of the conveyor may direct the puck along the direction of travel without leaving the conveyor. The puck may have movable parts for holding the object. The puck does not follow any track, but freely moves along a conveyor from station to station as desired. The puck may be gripped and/or aligned by a processing station (for example, a container filler) for a time, but the puck does not follow a defined track. Pucks can be problematical for various reasons. For example, because they are loose on a conveyor, pucks and their conveyed objects can be tipped over. Also, pucks may undesirably spin or move on a conveyor unless specifically guided since they do not follow any track. Further, different-sized pucks may be required if a line is to be converted from one object to another, which an be costly in terms of stocking and switching families of pucks."}
{"text": "Computing systems generally utilize non-volatile storage media, such as hard disk drives, to store information, and to ensure the availability of stored information across power cycles. In some instances, computing systems or applications executed by computing systems require that data be committed to a disk drive in a specific order, such that subsequent read requests to the disk drive result in retrieval of the expected data.\nHowever, the systems and protocols used to interact with storage devices such as disk drives may be unable to ensure such a specific order of writes. For example, multithreaded systems may allow multiple applications or processes to access a disk drive simultaneously or near-simultaneously, and in doing so, may reorder various operations performed by the applications or processes. Further, current disk drives or other non-volatile storage media can utilize volatile cache memory in an attempt to increase speed and efficiency of the drive. Use of cache memory, and protocols associated therewith, may result in reordering of operations prior to commitment to a non-volatile memory. Reordering of operations to the non-volatile memory may result in unexpected values being written to or read from the non-volatile memory by the computing system, and therefore result in errors.\nTo compensate for the out-of-order environment in which writes can be committed to the non-volatile memory, an application has the option to ensure that all of its previous writes have completed before completion of the next write by performing an application file operation flush, which causes all data of all previously completed writes to be forced out of the system and committed to the storage device's permanent media. However, implementation of a flush command on the storage device is generally time consuming, and can degrade performance. Consequently, many storage subsystems (or software/firmware/hardware associated therewith) decline to perform flush commands. Ignoring a program's (e.g., application's or operating system's) flush commands may increase performance relative to performing the commands, but results in potential data inconsistency.\nAn alternative technique for ensuring data consistency is to serialize all ordered writes throughout the host computing device and mark each one as a Forced Unit Access (FUA). Use of FUA writes ensures that each write is committed to non-volatile media before the command is completed. However, due to the delay in completion, FUA commands tend to be significantly slower than a typical write command. As a result, many storage subsystems (or software/firmware/hardware thereof) fail to respect a program's request to perform FUA writes."}
{"text": "The present invention relates to a radiographic image information reading method and a radiographic image reading apparatus, and in particular, to correction of an image to be read.\nRecently, there has been devised a method for obtaining radiographic image information without using a radiographic film composed of a silver halide light-sensitive material. According to one such method, radiation transmitted through a subject is absorbed in a phosphor of a certain kind, and then the phosphor is excited by light or heat energy, for example, so that radiation energy accumulated in the phosphor through the aforesaid absorption is emitted as a fluorescence which is detected to be an image. To be concrete, for example, U.S. Pat. No. 3,859,527 and Japanese Patent Publication Open to Public Inspection No. 12144/1980 (hereinafter referred to as Japanese Patent O.P.I. Publication) disclose methods employing a radiographic image conversion panel and showing a radiographic image conversion method with excitation light of visible light or infrared radiation wherein a radiographic image conversion panel in which a stimurable phosphor layer is formed on a support is used. Radiation transmitted through a subject is projected on the accelerating phosphor layer of the radiographic image conversion panel so that radiation energy corresponding to radiation transmittance on each part of the subject is accumulated to form a latent image. Then, the accelerating phosphor layer is scanned by the aforesaid exciting light so that radiation energy accumulated on each part of the panel is emitted to be converted to light, and then, its intensity is detected by a photoelectric transfer element such as a photomultiplier or a photodiode to obtain a radiographic image.\nFIG. 19 is an illustration showing how an image is recorded on a stimurable phosphor like that explained above. In the drawing, X-rays emitted from X-ray source 1 are narrowed down by diaphragm 2, and then are projected on subject 3. X-rays transmitted through the subject 3 enter accelerating phosphor 4 wherein a latent image of an image of the subject 3 is formed.\nFIG. 20 is a block diagram showing a structural example of a conventional radiographic image information reading apparatus which reads radiographic images recorded on an accelerating phosphor in the manner explained above. The numeral 101 is a semi-conductor laser light source for generating exciting light, and the semi-conductor laser light source 101 is driven by laser driver circuit 102 in the form of pulses in synchronization with an image clock signal from image clock generator 125. Laser beam LB generated from the semi-conductor laser light source 101 arrives at deflector 107 through monochromatic light filter 103, mirror 104, beam-forming optical system 105 and mirror 106. The deflector 107 is equipped with polygon mirror 109 driven by deflector driver 108, and it deflects the laser beam LB to cover a certain angle in the scanning area. The deflected laser beam LB is adjusted to be at a constant speed on a scanning line, and it scans on radiographic image conversion panel 111 employing a stimurable phosphor as a radiographic image information recording medium in the direction of arrow \"a\" through mirror 110. The radiographic image conversion panel 111 moves simultaneously in the sub-scanning direction (direction of arrow \"b\"), and thereby an overall surface can be scanned. Light emitted from the radiographic image conversion panel 111 after scanning by the laser beam LB are converged by converging means 112 and arrive, through filter 113 transmitting only a wavelength range of the light, at light detector 114 provided with a photomultiplier where the light are converted to analog electric signals (image signals).\nThe numeral 115 is a power supply which supplies high voltage to the light detector 114 (photomultiplier). An image signal outputted as an electric current from the light detector 114 passes through front end amplifier 116 to be voltage-amplified, and further passes through logarithmic amplifier 118 that converts radiation intensity signals into image density signals, filter 119 and sample hold circuit 120 that holds signals for a certain period in synchronization with image clock signals, and then is converted into a digital signal by A/D converter 121 to be sent to an external data processing apparatus through interface 124.\nIn the case of reading radiographic image information, however, there have been problems such as sensitivity irregularity of radiations and accelerating phosphors (irregularity in both main- and sub-scanning directions), irregularity of exciting light scanning system and light converging system (irregularity in main-scanning direction) and an influence of a change with age of a stimurable phosphor (irregularity in the sub-scanning direction). As a technology for overcoming the aforesaid problems, there has been a correction technology described in Japanese Patent O.P.I. Publication No. 234643/1985.\nHowever, in the technology described in the aforesaid patent, correction data for correcting the data read have been needed to be prepared for all pixels of a radiographic image conversion panel. Therefore, the capacity of a memory for the storage of correction data needs to be increased. When using a plurality of radiographic image conversion panels, in particular, memory capacity is increased in proportion to the number of panels, which is not preferable in practical use.\nFor solving the problems mentioned above, the present inventor proposed a novel method and a novel apparatus for reading radiographic image information (Japanese Patent O.P.I. Publication No. 158536/1988). According to this invention, correction data in a form of a line in the main-scanning direction and/or the sub-scanning direction are obtained from a solid image obtained through photographing without a subject and stored, and they are used for correcting image data in the case of photographing an image with a subject arranged. It is possible to save memory capacity greatly by having at least one list of correction data in the form of a line or a row without having correction data with total pixels two-dimensionally. Further, when obtaining correction data in the main-scanning direction or in the sub-scanning direction, it is possible to eliminate an influence of noise by averaging data of plural lines or rows.\nIn the invention, however, when streak-shaped noise (irregularity without position repeatability or an irregularity with poor reproducibility in terms of its position on an image) caused by vibration or irregularity of intervals on a polygon mirror appears on an image without a subject in the case of preparing correction data, the correction data can not be free from this influence. When image data is corrected by the use of correction data having the noise, the noise appears on the corrected image. As a cause of irregularity without position repeatability, there are irregularity in a light converging system and irregularity in an exciting light scanning system (defect and dust) for the main-scanning direction, and there are relative displacement (stationary vibration and shock from the outside) between a radiographic image conversion panel and a reading system, irregularity in an inclination angle of a polygon mirror and irregularity of reflection factor in mirror surfaces for the sub-scanning direction. With regard to such irregularity without repeatability, there have been some cases where the irregularity was increased when it was corrected.\nFIGS. 24(a-1) to 24(c-2) are illustrations of some problems in conventional technologies. In the drawing, f1 in FIG. 24(a-1) represents correction data, FIG. 24(a-2) is a frequency spectrum of f1, f2 in FIG. 24(b-1) represents image data, FIG. 24(b-2) is a frequency spectrum of f2, and f3 in FIG. 24(c-1) represents image data after correction and FIG. 24(c-2) is a frequency spectrum of f3. The axis of ordinate for f1-f3 represents a size of image data, the axis of abscissa represents a position, and the axis of ordinate for a frequency spectrum represents intensity, and the axis of abscissa represents a spatial frequency. An example shown in the drawing indicates an occasion wherein irregularity having position repeatability at a relatively gentle change and irregularity having no position repeatability appearing at a constant cycle at a relatively abrupt change are mixed. Since the irregularity having position repeatability appears similarly on correction data f1 and image data f2, effective correction is made on image data f3 after correction. On the other hand, irregularity having no position repeatability appears differently in terms of position between correction data f1 and image data f2. Therefore, effective correction can not be obtained, and irregularity is rather increased on the corrected image data f3. This is clear from the frequency spectrum shown in FIG. 24(c-2), wherein a peak of a specific frequency component is increased."}
{"text": "The present invention relates to a new and distinct variety of plum tree which will hereinafter be denominated varietally as \"Early Rosa\" plum tree and, more particularly, to a plum tree which produces clingstone fruit, which are mature for commercial harvesting and shipment approximately June 14 to June 19 in the northern San Joaquin Valley of California, having a sweet, rich flavor, moderately firm flesh and bright cherry red to dark purple skin coloration.\nThe \"Santa Rosa\" type plum has over a broad expanse of years developed a significant and seemingly permanent share of the fresh plum market. Due to its durability in the marketplace, there has been a well recognized and more or less continuous effort to develop new varieties of plum trees of the \"Santa Rosa\" type which bear fruit having enhanced, commercially appealing attributes including ripening dates which expand the market season for fresh \"Santa Rosa\" type plums beyond that which has previously been available. These efforts have not been consistently successful. The new variety of the present invention is believed to succeed admirably in these respects."}
{"text": "The invention pertains to a method for operating an exhaust gas aftertreatment with a diesel particulate filter and a device for controlling the exhaust gas aftertreatment. The invention also pertains to an exhaust gas aftertreatment and to an internal combustion engine.\nIt is known from the prior art that diesel particulate filters can be used to remove soot particles from an exhaust gas. Diesel particulate filters can comprise a fine-pored structure—e.g., a ceramic structure or, as described in US 2007-151,231 A, a fine-pored woven steel structure—on the walls of which the soot particles are deposited. It is known that diesel particulate filters can be regenerated. A distinction is made between passive regeneration and active regeneration; in the latter case, the soot particles are burned off at predetermined time intervals and/or after a predefinable trigger signal.\nIn an exhaust gas aftertreatment system with a passively regenerating diesel particulate filter, advantage is taken of the so-called CRT (Continuous Regeneration Trap) effect, and the diesel particulate filter is thus regenerated continuously, i.e., in particular without a fixed, predefined trigger signal. In the passive regeneration technique, the exhaust gas temperatures of the engine are sufficient under normal operating conditions to ensure the continuous removal of soot from the diesel particulate filter. Under certain conditions, e.g., special climate conditions or long, continuous periods of low-load operation, the regeneration process can be supported by more thorough measures. For example, the exhaust gas temperature of the engine and thus the soot burnoff can be increased significantly for a short time. US 2010-031,638 A, for example, describes implementing a passive regeneration by increasing the engine load. US 2011-265,456 A describes the possibility of increasing the soot burnoff by changing the combustion cycle and thus raising the temperature of the exhaust gas.\nTo determine a good time for regeneration, various mathematical and simulation or estimation methods for defining the soot load of a diesel particulate filter are described, some of which are quite complicated. It is known from WO 05/116413 that a neuronal network can be trained to determine the load state on the basis of the operating state of the engine, the differential pressure across the diesel particulate filter, and the exhaust gas values. The above-mentioned computational methods are comparatively complicated. In addition, there should be no need to carry out extra measurements to determine the state of the filter.\nIt is true that, to determine the load of the diesel particulate filter (DPF) and/or the best time for a regeneration process, the load state of the diesel particulate filter can be usefully determined by measuring the differential pressure (ΔP) across the diesel particulate filter. In the simplest case, regeneration by means of a thermomanagement measure as cited above could be initiated whenever a predetermined or calculated reference value for the differential pressure (ΔP) is exceeded.\nIt has been found, however, that the cause of the differential pressure is comparatively complex. The differential pressure increases over the service life of the diesel particulate filter not only as a result of the loading with soot but also as a result of a loading with ash; the accumulation of ash is a fundamental factor in determining the service life (T_L) of the diesel particulate filter. Whereas soot comprises essentially combustible carbon components, ash is defined as the incombustible component of a filter load, which is almost impossible to remove during operation or which at best requires considerable effort to do so. Thus, DE 1 002 951 describes measures for reducing the amount of ash.\nDuring operation, however, the increase in the differential pressure caused by the ash load leads in particular to the situation that, either a reference value pertaining to the differential pressure for initiating regeneration is reached at an increasingly earlier point, or a reference value calculated on the basis of the combustion ends up being too low. As a result, this leads in the case of a passive regeneration system to the initiation of thermomanagement more frequently than would in fact be necessary, i.e., it would lead regularly to the initiation of thermo-management measure even though the soot load of the diesel particulate filter is not yet critical. It is known from DE 12 034 340 A1, for example, that a combustion calculation for a known fuel specification can be carried out, which also takes into account the ash residues in the diesel particulate filter. It is also desirable, however, to have available an improved method and a device for operating an exhaust gas aftertreatment system by means of which the influence of ash on the differential pressure across a diesel particulate filter of the exhaust gas aftertreatment system can be taken into account, i.e., in particular so that the influence can be determined, in the case of passively regenerating diesel particulate filters."}
{"text": "1. Field of the Invention\nThe present invention relates to a photoelectric converter device in which the wavelength characteristic is adjusted by means of a wavelength filter, and to a method for manufacturing the photoelectric converter device.\n2. Description of the Related Art\nIn recent years, a chip size package (CSP) comprising wiring lines drawn out from a side surface of an element is widely used to miniaturize devices such as photoelectric converter devices which include photoelectric converter elements.\nFIGS. 12A and 12B illustrate top and bottom external views, respectively, of a semiconductor integrated device configured as a CSP. The CSP is constituted by adhering an upper support member 14 and a lower support member 16 on upper and lower surfaces, respectively, of a semiconductor chip 10 by means of resin layers 12 composed of epoxy or the like. External wiring lines 18 are drawn out from side surfaces of the semiconductor chip 10 sealed between the upper support member 14 and the lower support member 16. The external wiring lines 18 are connected to ball-shaped terminals 20 provided on the bottom side of the CSP.\nFIG. 13 shows a cross-sectional view of a conventional photoelectric converter device including photoelectric converter elements configured as a CSP. The semiconductor chip 10 is formed as a lamination of a semiconductor substrate 10a having photoelectric converter elements formed on its upper surface, color filters 10b and 10c disposed to at least partially cover the photoelectric converter elements for absorbing or reflecting light in predetermined wavelength regions, and a planarization film 10d for planarizing surface unevenness. The upper support member 14 is adhered to the upper side of the semiconductor chip 10 via a resin layer 12, while the lower support member 16 is adhered to the lower side of the semiconductor chip 10 via another resin layer 12. With this arrangement, by employing a light-transmitting substrate made of glass or the like as the upper support member 14, external light can be introduced into the photoelectric converter elements formed on the upper side of the semiconductor chip 10.\nIn the above-described arrangement, each of the color filters, which is a visible light filter for permitting transmission of visible light, is composed of a filter material which mainly permits transmission of red (R), green (G), or blue (B) wavelength region. FIG. 14 illustrates sensitivity characteristics of a photoelectric converter element obtained when using red (R), green (G), and blue (B) color filters that are typically employed. A typical color filter exhibits a relatively high transmittance in the infrared region exceeding 700 nm in addition to in a wavelength region of the corresponding color. Further, a silicon photoelectric converter element shows responses in the infrared region exceeding 700 nm. In FIG. 14, wavelengths of incident light into the filters are given on the horizontal axis, while sensitivity levels obtained at various wavelengths are given on the vertical axis. In FIG. 14, line A denotes a characteristic obtained when using a color filter corresponding to the red (R) wavelength region, line B denotes a characteristic obtained when using a color filter corresponding to the green (G) wavelength region, and line C denotes a characteristic obtained when using a color filter corresponding to the blue (B) wavelength region. Furthermore, line D denotes a typical example sensitivity characteristic, with respect to various wavelengths, of a photoelectric converter element (without color filters) formed on a silicon substrate.\nAs can be seen, when a color image is captured using these color filters, the obtained output signals for the respective colors are influenced by noise generated due to infrared components included in the incident light. Consequently, when those output signals are combined to produce an image, correct white balance of the image cannot be achieved.\nIn order to avoid the above-described problem, a structure as shown in FIG. 15 is conventionally used. In FIG. 15, an infrared cut-off filter 102 is arranged so as to cover the incident surface of a photoelectric converter device 100 configured as a CSP, while a lens 104 is provided to direct incident light into the photoelectric converter elements.\nAs can be seen in FIG. 15, the structure including the infrared cut-off filter 102 provided over the photoelectric converter device 100 is disadvantageous in that the overall device size becomes very large.\nAnother disadvantage is that the cost of the infrared cut-off filter 102 is very high because the filter 102 is fabricated by laminating many layers of infrared-absorbing materials having different characteristics on a glass substrate. Use of the infrared cut-off filter 102 is one significant factor which increases the fabrication costs of the photoelectric converter device 100. Furthermore, when employing a configuration in which the infrared-absorbing films are formed on the surface of the glass substrate, the infrared cut-off filter 102 often becomes damaged due to handling during the fabrication process. Moreover, the difference in contraction ratio between the glass substrate and the infrared-absorbing films may result in curving and partial peeling of the infrared cut-off filter 102."}
{"text": "1. Field of the Invention\nEmbodiments of the invention are related to tools used for testing application software. More specifically, embodiments of the invention provide a coverage tool for analyzing database-aware applications.\n2. Description of the Related Art\nCurrently, coverage tools exist for testing software applications. Typical coverage tools analyze the coverage of the program source code of the application during a testing process. That is, they determine which parts of the program source code (e.g., classes, methods, routines, functions, etc.) are used, and are therefore tested, during a given test.\nHowever, current coverage analysis tools do not evaluate the coverage of parts of a database (e.g., tables, columns, rows) that may be accessed while testing the application. For example, if a high percentage of a project's code is being tested, but a high percentage of the database query strings are in the untested source code, the coverage analysis may be misleading. Similarly, coverage analysis does not include analysis of what tables, rows, elements, or attributes are accessed by the database query strings, because the coverage remains unaware of the semantic content of the strings.\nAccordingly, as the foregoing discussion illustrates, there remains a need for a coverage analysis tool for testing database-aware applications."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of eye drop dispensing devices and more specifically to an eye drop bottle holder with resilient legs.\n2. Description of the Related Art\nDevices for dispensing eye drop solutions are known. Generally, a bottle of eye drop solution includes a drop dispenser that is built into the exit orifice of the container. To dispense the solution, the user squeezes the bottle forcing solution out of the exit orifice and into his or her eye. Many users have trouble with dispensing eye drops from standard dispensing bottles. The user has a tendency to blink when the drop is about to enter the eye, causing the drop to miss the eye and land on a closed lid. Therefore, eye drop solution is frequently wasted due to the user blinking during the attempted application and the user ends up with eye drop solution streaming down his or her face.\nA number of efforts have attempted to resolve the above mentioned problem. Thomas Keen, in his U.S. Pat. No. 4,543,096, discloses a dispenser with an eyelid opening device. The user is required to place a pair of lid spreading arms dangerously close to his eye and then to press a lever arm to keep the eyelids apart. Thomas Sherman, in his U.S. Pat. No. 6,371,945, discloses an attachment for a bottle that includes a ring intended to help align the bottle with the eye. However, no attempt is made to hold the eyelid open. Gary Campagna, in his U.S. Pat. No. 3,934,590, shows a tripod like device for aligning the solution bottle over the user's eye. No attempt is made to hold the lid open. James Davidian, in his U.S. Pat. No. 6,595,970, shows a device for dispensing eyed drops. He proposes a dispensing arm, one side of which includes an indentation that receives the user's nose, the other side of which accepts a dispensing bottle. The bottle includes a pair of arms which, when squeezed, impinge on the side walls of the bottle forcing solution out of the bottle and into the user's eye. No attempt is made to hold the user's eyelid open. The U.S. Pat. No. 4,692,834 of Julia Clifford et al shows a dispenser that facilitates the user's amount of drops that exit a solution holding bottle. This bottle has retractable apertures that capture and release a drop of solution. James Walsh, in his U.S. Pat. No. 4,927,062, U.S. Pat. Nos. 6,041,978 and 6,010,488 and U.S. Pat. No. 4,834,727 as well as U.S. Pat. No. 5,902,292, all attempt to position an eye drop bottle in a correct location above a person's eye, but none of them includes a means to help hold the user's eye lids in the spread apart, open position. U.S. Pat. No. 4,321,916 discloses an eyelid retractor that is used during ocular surgery or the like. It is not designed to be used by a user of eye drop solution.\nNone of the above cited inventions safely holds the user's eyelids open while dispensing eye drops from a standard eye drop bottle. Additionally, none of the above mentioned patents allows the user to dispense a portioned amount of eye drop solution in an automatic and repeatable fashion. None of the above cited inventions dispenses a precise amount of eye drop solution and simultaneously holds the user's eyelids open while doing so."}
{"text": "Packing material is widely used to prevent fluid leakage around an operating member in a housing with fluid, such as a rotary shaft or a sliding stem in fluid control valves or in a reciprocating pump member. Normally such packing is formed of a resilient material and is placed under a static load by being bolted into position within a packing box around the operating member. In other instances the packing is subjected to spring loading in what is known as a live loaded packing configuration. Live loaded packing is particularly useful in attempting to prevent leakage of undesired fluids into the environment. Also, it is desired to use packing material formed of polytetrafluorethylene (PTFE) because of its low frictional impact on an operating member.\nGenerally, because of its low elasticity, it has been found necessary to continually load PTFE packing in the containment assembly in order to enable the packing to set in and maintain a reliable and tight fluid seal against the operating member and the housing. However, in attempts to use PTFE packing in a live loaded packing configuration at extremely elevated temperatures and pressure conditions, the PTFE packing tends to deform and cold flow and eventually extrude out of the sealing area. Substantial extrusion of the PTFE packing leads to fluid leakage and increased wear and possible damage to the operating member and packing.\nFor example, fluid control valves are operated in environments where the operating pressure is less than about 1000 psi (6895 kPa) and the temperature of the fluid is less than about 200.degree. F. (94.degree. C.). However, certain applications of fluid valves require use in much higher pressure and temperature operational environments such as greater than 1000 psi (6895 kPa) and around 450.degree. F. (232.degree. C.). As an example, in the control of boiler feed water in power generating plants there is a requirement for fluid valves to reliably operate in pressures around 7200 psi (49644 kPa) and in temperature ranges of around 450.degree. F. (232.degree. C.).\nNormally, valves operating at such pressure and temperature extremes use graphite packing material. (See for example U.S. Pat. No. 4,364,542 suggesting the use of a graphite packing material at temperatures above 250.degree. F. (107.degree. C.). However, the use of graphite packing leads to increased valve stem friction as the valve stem strokes through the packing thereby severely limiting valve operation and leading to increased wear on the packing material.\nWhile desired, prior attempts to utilize PTFE packing at elevated pressure and temperature ranges was normally not recommended nor found feasible, particularly under live loaded packing conditions as indicated previously. In such attempts, it was found that the PTFE packing material extruded due to the severe pressures and movement of the operating member such as a valve stem, which tended to remove the material from the packing in long strands of string-like extrusions. Eventually, the continued removal and extruding of material from the packing led to a loss of sealing ability for the packing and thereby requiring early replacement of the packing material and possibly the operating member itself.\nIn U.S. Pat. No. 4,886,241, assigned to the same assignee as herein, there is described a packing containment structure for use in fluid valves at high pressure and high temperature conditions around 7200 psi (49644 kPa) and 450.degree. F. (232.degree. C.). This patent describes the use of PTFE packing material with an anti-extrusion wiper ring on each side of the packing material. The wiper rings are formed of high temperature organic or inorganic fiber material with a nitrile elastomeric binder. Such a packing containment structure has been found suitable for use at the described high pressures and temperatures where the fluid is water or other non-reactive fluid. However, due to the fiber material used for the anti-extrusion wiper rings, this packing containment structure is not recommended for use where the fluid is chemically active and where extremely stringent fluid control requirements must be met.\nBecause of the highly inert properties of PTFE it is desired to use such material for the anti-extrusion wiper ring in the configuration of the packing containment described in the aforementioned U.S. Pat. No. 4,886,241. However, the known undesired deformation and cold flow properties of PTFE negatively impact on the use of this material at the extreme temperatures and pressure conditions presently under consideration. Furthermore, while the aforementioned U.S. Pat. No. 4,886,241 describes how to use desired PTFE packing material at these elevated temperature and pressure conditions, only a fiber material with a nitrile elastomeric binder, i.e., a non-PTFE material, was suggested for the anti-extrusion wiper rings.\nAccordingly, it is desired to provide a packing containment for use with chemically reactive fluids in high pressure and high temperature conditions such as greater than 1000 psi (6895 kPa) and around 450.degree. F. (232.degree. C.) with PTFE packing material and which inhibits the extrusion of PTFE."}
{"text": "1. Field of the Invention\nThe present invention relates to apparatuses and processes for mixing, coating and drying a plurality of objects. Most preferably, this invention relates to apparatuses and processes for distributing and drying a sugar coating to the surfaces of confectionery.\n2. Related Background\nMany edible items (xe2x80x9cpiecesxe2x80x9d) have an outer coating comprised of sugar and other edible material. Such outer coatings, or candy shells, serve many functions, including preserving the center confectionery, as well as supplying an appealing look and taste to the edible item. Candy shells also give an identified xe2x80x9ccrunchxe2x80x9d when eaten.\nMany devices are known for applying a coating to edible pieces. For example, U.S. Pat. No. 5,507,868 discloses a typical prior art sugar coating apparatus. The reference describes a device having a rotary drum panner having a plurality of air holes and hollow shaft portions. Drying gas is supplied through one of the hollow shaft portions and directed downward toward an accumulated layer of particles for drying.\nU.S. Pat. No. 4,050,406 discloses a sugar coating apparatus which uses an annular mesh trough for containing objects for coating. The apparatus is mounted on springs and is vibrated to move the objects around the annular container. Drying air is supplied through the center of the apparatus directly to the mesh. The airflow then flows straight through the mesh to exit the top of the apparatus. The airflow is not a tangential flow.\nAnother prior art device is disclosed in U.S. Pat. No. 5,158,804. This patent discloses a coating apparatus and process for coating small quantities of tablets (drugs during the initial stages of formulation and development). A spray coating is applied from above the tablets on a mesh and drying air is supplied upwardly from below the mesh. In addition, the mesh is vibrated to make the tablets bounce up and down.\nStill other prior art devices fluidize objects for coating. For example, in U.S. Pat. Nos. 5,792,507 and 5,296,265, disclose coating apparatuses and processes which use a rotating disc in conjunction with an air flow at the periphery of the disc near the inner wall of a cylindrical container to fluidize the particles to be coated.\nThe above apparatuses and processes are insufficient for economically coating a multitude of confectionery since the costs associated with building, operating and maintaining them are relatively high. In addition, the prior art processes require a relatively long period of time to coat confectionery, which results in the expenditure of a greater amount of energy per unit mass of coated material. Further longer coating times reduce productivity of a piece of equipment resulting in higher cost per unit mass of product.\nThus, there is a need for an efficient, compact, and faster apparatus, which consumes less energy than existing apparatuses, and which produces coated confectionery in shorter period of time.\nThe present invention addresses the above concerns and presents new and novel apparatuses and processes for coating and drying a plurality of particles. Moreover, the present invention lends itself to any process which benefits from good mixing of particles within a deep bed of product, especially for coating, mixing and drying a hard-panned sugar shell applied to confectionery.\nThe factory floor area required for the present invention is greatly reduced over existing apparatuses, yielding additional cost savings.\nIn one aspect of the present invention, an apparatus for heating and/or cooling a plurality of particles includes a generally circular bed for holding the particles, the bed including a substantially centered annular inner bed wall, an outer bed wall and an irregular bottom. The irregular bottom is an uneven surface which may or may not be perforated and most is preferably a screen. The apparatus also includes a drive mechanism rotatably affixed to the bed, where rotation of the drive mechanism at a predetermined rate produces a displacement of the bed horizontally and vertically resulting in a particle flow in the x, y and z directions. The particle flow in the x and y planes for a majority of the particles is in a dominant direction. The apparatus also includes an air plant which provides a heated and/or cooled airflow tangentially below the screen. The airflow moves through the screen, and into the bed of particles.\nThe air plant may include a pump for providing an airflow into a duct, a dehumidifier in fluid communication with the pump, a heater positioned downstream from the dehumidifier and in-line with the duct, and a generally circular intake manifold having an outer manifold wall and an airflow inlet in fluid communication with the duct and tangentially arranged on the outer manifold wall. The intake manifold is positioned below the bed and includes an outer circumference substantially corresponding to the bed.\nAn outlet may be included along the outer wall of the bed and be removably covered by a tangentially positioned cover provided around a portion of an inner surface of the outer wall. The cover includes a first end affixed to the inner circumference of the outer wall downstream of the outlet in the dominant direction, and a second movable end adjacent the inner surface of the outer wall when the outlet is closed. The movable end is be positioned upstream from the outlet in the dominant direction and is opened when the movable end is moved toward the center of the bed. The cover and outer wall form a funnel having an entry for accepting particles from the dominant direction.\nIn another aspect of the present invention, a process for heating or cooling a plurality of particles in an apparatus as described in the previous aspect includes the step of displacing the bed at a predetermined frequency vertically and horizontally to produce a particle flow in the x, y and z directions. Particle flow in the x and y planes for a majority of the particles is in a dominant direction. The process also includes the step of heating and/or cooling an airflow to a predetermined temperature and directing the airflow tangentially to the bed via a generally circular intake manifold substantially corresponding in circumference to the circular bed and positioned below the screen. The airflow flows in a circumferential pattern around the interior of the intake manifold, through the screen and into the bed.\nIn yet another aspect of the present invention a process for coating and drying a plurality of particles in the apparatus, as described in the first aspect, includes the step of displacing the bed as described in the previous aspect to produce a similar particle flow. Other steps of the process include, in the case of coating particles, applying a coating material from a nozzle positioned above the bed, the application of the coating occurring during a first predetermined period of time, providing a temperature controlled airflow and directing the temperature controlled airflow tangentially to the bed via a generally circular intake manifold which substantially corresponds in circumference to the circular bed and is positioned below the irregular bottom. The airflow being provided after a second period of time has passed after the coating material has been applied. The airflow moves in a circumferential pattern around the interior of the intake manifold, through the irregular bottom and into the bed.\nWhile a preferred embodiment provides the airflow tangentially to the bed in a direction opposite to the dominant mixing direction, air flow in the direction of the dominant mixing direction may also be used.\nThe high degree of mixing of the individual particles in the x, y and z directions in the present invention ensures consistent treatment of all particles with the desired process. This high mixing reduces piece to piece variation which that results when individual pieces remain in localized areas within the process volume when drying, coating or other conditions differ from the average.\nAccordingly, it is a feature of the present invention that a high degree of shear and contact is created between individual pieces of the confectionery product as a result of the particle flow created. This results in faster coating and drying times since confectionery are mixed more thoroughly and evenly and are more uniformly coated, without damaging the confectionery or causing rough surfaced candies.\nThe tangential air injection into the cylindrical area creates a centrifugal pressure gradient which encourages the majority of the air flow to occur through the outer area of the bed where the majority of the particles are present. This maximizes the desired transfers between the air and the particles.\nThe present invention optimizes airflow and temperature to maximize evaporation rate. This results in a cycle time for some confectionery which is 20-40% of the time of traditional apparatuses and methods, even with relatively high air humidity dew-points approximately between 4-12xc2x0 C. By carefully monitoring and controlling the temperature of the product in the bed throughout the coating cycle, a hotter supply of air on average can be used, thereby increasing the amount of moisture that can be removed per cfm.\nThus, it is a further feature of the present invention that the dehumidified and heated airflow provided tangentially to the bed provides maximum evaporative effect to the particles within the bed so that they may be dried faster than other apparatuses and processes disclosed in the prior art.\nThe speed of coating and drying occur so much faster than in existing techniques, that lower-solids coating solution may be used. High-solids coatings, used in existing techniques, are generally used so that less moisture will be required to be removed from the coated particles after application. High-solids coatings, however, must be kept at heated higher temperatures. Thus, use of lower-solids coating solutions eliminates the need for keeping coatings hot, thereby eliminating the need to heat the coatings and also eliminates the need for water jacketing of piping and containers for ferrying and storing the hot coating. The elimination of hot coating also contributes to significant energy savings due to less heat loss to the environment air."}
{"text": "Current interaction with televisions, television set-top boxes, and/or other television devices generally involve a viewer controlling the devices with various remote controls. A typical home may include several televisions and/or television set-top boxes located throughout various rooms of the house. In addition, several different remote control devices are needed for all of the different brands and locations of the television devices. Typically, a user wanting to watch television enters a room of the house and needs to find the remote controls that correspond to the television and/or television set-top box that are located in the particular room of the house."}
{"text": "Training simulators are frequently used to provide realistic training of complex skills. One advantage of a training simulator over real equipment is that it can provide training more cost-effectively. Since visual information is often necessary to guide learned behaviour, an interactive visual system is needed to supply environmental cues. However, the visual system of a training simulator is usually its most expensive component. Low cost solutions are thus needed to improve cost-effectiveness.\nAn alternative to systems currently in widespread use is an interactive, videodisc-based visual system which can provide realistic backgrounds and afford a limited amount of dynamic interaction with previously recorded scenes. In order to allow effective training of visual-motor skills, it is necessary that the visual system be capable of responding to this interaction in real-time. Also, it is sometimes necessary to provide additional, dynamic visual information to supplement recorded training materials. This is required because it is not always possible to anticipate fully the temporal and spatial occurrence of a visual event, or its exact visual characteristics.\nOne solution is to combine low-cost imagery generated by computer with video recordings. When necessary, the computer generated imagery is superimposed upon, or keyed over, the prerecorded background. A system of this type is appropriate to the design of training simulators meant to assist the learning of tasks which rely on visual depictions, the area in which visual changes occur being limited and not requiring detailed imagery to be produced continuously. There is a fairly large range of military tasks in need of training which meet these requirements. For example, in training tank gunnery the target is the focus of attention and little computer-generated imagery needs to be used to simulate the trajectory and fall of shot. However, it is often very expensive to train personnel for these tasks, and accordingly the system is particularly appropriate for such training purposes. Such a system could also easily be adapted for entertainment purposes, in a wide variety of arcade games."}
{"text": "Leakage power dissipation is an important issue in the design of large scale integrated circuits such as microprocessors. Increased leakage, e.g., as a result of decreasing transistor dimensions, has resulted in the need for increased keeper sizes in dynamic circuits. Unfortunately, increased keeper sizes are not necessarily needed for all chips at all times in the same manner and can have undesirable results such as higher contention, higher power consumption, higher cross-over (short-circuit) current, and diminished performance. Accordingly, an improved keeper solution is desired."}
{"text": "The use of various steroids to increase the intensity of the heartbeat is known. A number of steroid glycosides have been found useful in cardiotherapy, but the corresponding aglycons (genins), which share the characteristic steroid structure of the glycosides are convulsive poisons.\nFor purposes of definition, a cardiotonic glycoside is one which increases the intensity of the heart beat but may decrease the rate of heartbeat. These glycosides may be designated by the following general formula: ##SPC1##\nThe corresponding aglycons have a hydroxy group at position 3 rather than a glycoside conjugate.\nA wide number of sugars may be employed to form cardiotonic glycosides, as long as the steroid nucleus is a cardenolide which contains a five membered lactone ring at the C.sub.17 position. Particular cardiotonic glycosides useful in the treating of heart disceases and which are readily available include digoxin, digitoxin and gitoxin. These compounds differ at the C.sub.12 and C.sub.16 position of the steroid (R.sub.1 and R.sub.2), in that digitoxin contains a hydrogen substituent at each position, digoxin contains a hydroxyl at C.sub.12 and a hydrogen at C.sub.16 while gitoxin contains a hydrogen at C.sub.12 and hydroxyl at C.sub.16. The glycoside portion of each of these compounds comprises three digitoxose sugars. The 3', 3\", 3'\" and 4'\" positions bear hydroxyl groups while the 5', 5\" and 5'\" bear methyl substituents. Other sugars which may be found in cardiotonic glycosides include the lanatoside series which differs from the digitoxose in that the 3'\" position bears an acetyl while the 4'\" position contains an additional glucose sugar. When this glycoside is conjugated to digoxigenin it is designated lanatoside A, while if it is conjugated with gitoxigenin it is designated lanatoside B, and if it is conjugated with digitoxigenin it is designated lanatoside C.\nOther sugar units which may give a compound which has cardiotonic activity, when they are conjugated to a cardenolide steroid, include rhamnose, antiarose, digitalose, thevetose, talomethylose, cymarose, oleandrose, sarmentose, boivinose and diginose. By combining three of these units with an aglycon, various cardiotonic glycosides may be formed. The aglycons alone are usually convulsive poisons not useful in medicine. For a more thorough discussion of cardiotonic glycosides and aglycons, see \"Rodd's Chemistry of Carbon Compounds\", Second Edition, Vol. 2, part D, chapter 17, and \"Chemistry and Metabolism of Digitalis\", Digitalis, edited by Charles Fisch and Borys Furawicz, Grune & Stratton, Inc., New York, 1969.\nAs previously mentioned, although the glycosides have been found to exhibit much less toxicity than the corresponding aglycon, in many instances the difference between a therapeutic and a toxic dose of a glycoside is only in the order of a few micrograms per kilogram of blood. Digitalis intoxication occurs frequently due to the small differences between therapeutic and toxic quantities. The therapeutic quantity of any cardiotonic glycoside cannot be predicted in advance of medication since each patient may react differently to the dosage, and retain different quantities of the drug in the blood stream.\nDigitalis intoxication occurs with increasing frequency in elderly heart patients, since diseased renal function in the aged can result in higher and sometimes toxic levels of cardiotonic glycoside residue in the blood. Also, there may be some overlapping between therapeutic and toxic blood levels due to susceptability of the myocardium to digitalis. For example, the patient whose myocardium is diseased may show signs of toxicity despite therapeutic blood levels. Other clinical conditions contributing to digitalis intoxication are coronary ischemia and potassium depletion. The problems of digitalis intoxication are compounded in that its symptoms may often be mistaken for fatigue or restlessness.\nGenerally, accepted therapeutic levels for digoxin range from 0.8 to 2.4 ng/ml while toxic levels range from 2.1 to 8.7 ng/ml Toxic versus non-toxic levels of digitoxin also overlap, i.e., 3.0 to 39.0 ng/ml may be therapeutic while 26.0 to 43.0 ng/ml may be toxic. Higher levels of digoxin and digitoxin are generally lethal.\nRadioimmunoassay for cardiotonic glycosides has proved extremely effective in detecting undermedicated, adequately medicated and over-medicated or toxic levels of the drug in blood serum. That sampling technique comprises the use of a radioactive steroid glycoside or aglycon and an antibody which will bind both the radioactive steroid and the cardiotonic glycoside in the blood on a competitive basis. By counting the amount of radioactive material reacting with the antibody, the amount of cardiotonic glycoside in the blood may be calculated. From the results, it can be quickly determined whether the patient has been undermedicated, properly medicated or over-medicated and is suffering digoxin toxicity.\nFor radio labelling of steroids, 125-Iodine or 131-Iodine are preferred, since they exhibit high specific activity, while .sup.14 C and .sup.3 H are also available. The use of compounds exhibiting a high specific activity permits the use of considerably less material for the same counting efficiency. Counting with such isotopes can be achieved by liquid scintillation or gamma ray spectroscopy counting procedures.\nThe sampling procedure preferably consists of the use of an antibody polymer in tablet form, which antibody has been raised in an animal host by the administration of an appropriate antigen. The radioactive steroid and the cardiotonic glycoside isolated in the blood serum then complete with each other for a limited number of binding sites on the antibody to form an insoluble complex. The antigen/antibody complex is then separated via centrifuge techniques and the radioactive portion is measured, as for example, by gamma counter. From the difference in radioactivity in the test specimen, compared to standards of known cardiotonic glycoside concentrations in blood, the level of cardiotonic glycoside in the test sample may be accurately computed.\nFor a more complete discussion of radioimmunoassay principles, see A. R. Midgday, Jr. and G. D. Niswender, in \"Karolinska Symposia on Research Methods in Endocrinology, 2nd Symposium: Steroid Assay by Protein Binding,\" The Reproductive Endocrinology Research Unit, Stockholm (1970), 320-333 and W. D. Odell and W. H. Daughadey, \"Principles of Competitive Protein-Binding Assays,\" J. B. Lippincott Company, Philadelphia (1971). For a complete discussion on the preparation of protein-digoxin conjugates necessary for the production of antibodies, see \"Digoxin-Specific Antibodies\", Butler et al., Proc. Natl. Acad. Sci. U.S. 57:71.\nThe method most commonly used in conjugating a radioactive isotope, such as 125-Iodine, to a steroid is via an electrophilic substitution reaction. Because of the high susceptability of phenol derivatives to electrophilic substitution, since the hydroxyl is an activating group, radioactive steroids have been synthesized by conjugating tyrosine or tyrosine methyl ester (TME) onto the steroid, after which the compound may be reacted with the radioactive isotope (.sup.125 I).\nOliver et al., Journal of Clinical Investigation, 47, 1035 (1968) teach the synthesis of a digitoxigenin-3-0-succinyl-iodine-125 tyrosine methyl ester for use in cardiotonic glycoside radioimmunoassay procedures. Such a compound, although suitable for measuring the level of the drug in blood serum, suffers from the disadvantage that it is extremely difficult and expensive to produce. The succinylation step of the C.sub.3 -hydroxyl group of the genin, which is used at the binding site for the TME and iodine-125, proceeds very slowly and may require reaction time on the order of months. It is believed that the slow reaction time associated with the succinylation of the C.sub.3 -hydroxyl, is due to the fact the C.sub.3 -hydroxyl group is fixed in an axial position and offers a large amount of steric hindrance. Additionally, although such a reaction procedure is suitable for the synthesis of digitoxigenin derivatives, where the C.sub.12 substituent is a hydrogen, it is not per se suitable for the synthesis of radioactive digoxigenin derivatives. Digoxigenin has at the C.sub.12 position and equatorial hydroxyl substituent which will more readily be succinylated than the axial-3-hydroxyl group, since it offers less steric hindrance. Thus succinylation of the 3-hydroxyl of the genin not only proceeds slowly but for many compounds produces little or no yield of a useful immunoreactive derivative unless other positions are blocked. In addition, Stall et al., \"The Specificity of the Digoxin Radioimmunoassay Procedure\", Res. Commun. Chem. Pathol Pharmacol, 4, 503 (1972) shows that digoxigenin has a lower affinity for antibody than digoxin itself. They show a crossreactivity of about 34% wherein the optimum ratio of binding ability of radiolabelled to unlabelled antigen should be 1:1. Thus it would be expected that a radioactive glycoside will bind on a more competitive basis with an antibody than the corresponding aglycon.\nAlthough Draws et al., \"Faster and Easier Radioimmunoassay of Digoxin,\" Clinical Chemistry, Vol. 20, No. 3, 343-347 (1974) discusses the use of a .sup.125 I-tyrosine methyl ester of digoxin, there is no indication that such a compound should be a sugar hydroxyl substituted material or any indication of how such a compound might be synthesized.\nIt is an object of this invention to provide a process of synthesizing a radioactive steroid which will complete on a predictable basis for sites on an antibody with cardiotonic glycosides to be measured.\nIt is another object of this invention to provide a process of synthesizing such radioimmunoassay compounds via an improved and simplified procedure.\nIt is yet another object of this invention to provide a radioimmunoassay compound having a high specific activity and improved stability upon storage."}
{"text": "In most later-developed oilfields, a plethora of information may be available to the production engineer to assist in making decisions, such as deciding which wells should be scheduled for workover operations. In many cases, the amount of information may cloud the decision making process. For example, when comparing just two wells, information such as water cut, oil production rate, gas production rate, and lift parameters as between the two wells may give mixed signals as the ranking the two wells for workover operations. The problem compounds exponentially when considering a significant number of wells and/or patterns of wells."}
{"text": "1. Field of the Invention\nEmbodiments of the invention relate to a computer system; and more specifically, to the use of byte-addressable non-volatile random access memory for BIOS storage.\n2. Description of the Related Art\nA. Current Memory and Storage Configurations\nOne of the limiting factors for computer innovation today is memory and storage technology. In conventional computer systems, system memory is typically implemented by dynamic random access memory (DRAM). DRAM-based memory consumes power even when no memory reads or writes occur because it must constantly recharge internal capacitors. DRAM-based memory is volatile, which means data stored in DRAM memory is lost once the power is removed.\nWith respect to mass storage, conventional mass storage devices typically include non-volatile magnetic media (e.g., hard disk drives) and/or flash memory (also referred to as “flash”) (e.g., solid state drives (SSDs)). These storage devices are block-addressable, which means that a single byte of storage cannot be accessed individually. Rather, bytes are read and written as multi-byte (e.g., 16 bytes) blocks of data. Generally, these storage devices are considered I/O devices because they are accessed by the processor through various I/O adapters that implement various I/O protocols. These I/O adapters and I/O protocols consume a significant amount of power and can have a significant impact on the die area and the form factor of the platform. Additionally, for portable or mobile devices (e.g., tablet computers, cameras and mobile phones) that have limited battery life, their storage devices (e.g., Embedded Multimedia Card (eMMC) and Secure Digital (SD) card) are typically coupled to the processor via low-power interconnects and I/O controllers in order to meet active and idle power budgets. These interconnects and I/O controllers cannot consistently deliver the bandwidth required for a satisfying user experience.\nWith respect to firmware memory, a conventional computer system typically uses flash memory devices to store persistent system information that is read often but seldom (or never) written to. For example, Basic Input and Output System (BIOS) images are typically stored in a flash memory device. Flash memory devices that are currently available in the market generally have limited speed (e.g., 50 MHz). This speed is further reduced by the overhead for read protocols (e.g., 2.5 MHz). In order to speed up the BIOS execution speed, conventional processors generally cache a portion of BIOS code during the Pre-Extensible Firmware Interface (PEI) phase of the boot process. However, the processor cache has a very limited amount of capacity. Thus, the amount of BIOS code that can be used for the initial system configuration is also very limited. The size limitation of the processor cache places a significant restriction on the size of the BIOS code used in the PEI phase (also known as the “PEI BIOS code”). Therefore, the PEI BIOS code cannot be easily extended to support a large mix of memory configurations and multiple processor families. With the rising demand for initialization of processors, process interconnects, memory and storage implemented with various technologies and multiple processor families, the need for more versatile PEI BIOS code also grows. One solution is to build a larger processor cache for caching the code. However, the size of the processor cache cannot be easily increased without a negative impact on the rest of the system.\nB. Phase-Change Memory (PCM) and Related Technologies\nPhase-change memory (PCM), also sometimes referred to as PCME, PRAM, PCRAM, Ovonic Unified Memory, Chalcogenide RAM and C-RAM, is a type of non-volatile computer memory which exploits the unique behavior of chalcogenide glass. As a result of heat produced by the passage of an electric current, this material can be switched between two states: crystalline and amorphous. Recent versions of PCM can achieve two additional distinct states, effectively doubling memory storage capacity. PCM is one of a number of new memory technologies competing in the non-volatile role with flash memory. Flash memory has a number of practical problems which these replacements hope to address.\nFor example, PCM can offer much higher performance in applications where writing quickly is important, in part because the memory element can be switched more quickly, and also because individual bits may be changed to either 1 or 0 without the need to first erase an entire block of cells (as is the case with flash). The high performance of PCM makes it potentially very beneficial in non-volatile memory roles that are currently performance-limited by memory access timing.\nAdditionally, while PCM devices degrade with use (like flash), they degrade much more slowly. A PCM device may survive approximately 100 million write cycles. PCM lifetime is limited by mechanisms such as degradation due to GeSbTe (GST) thermal expansion during programming, metal (and other material) migration, and other mechanisms."}
{"text": "This invention relates to a radar system which is defined for the purpose of this specification as apparatus for deriving information about the position or movement of a feature in a medium surrounding the apparatus by transmitting energy, through the medium, to a feature in said medium and observing the energy travelling to the said apparatus from the aforesaid feature. Whilst the invention is particularly applicable to apparatus which transmits and receives microwave radiation it is also applicable to acoustic or optical systems. Also, whilst the invention is particularly applicable to radar systems which sense the direction and time of arrival of the energy it is also applicable to systems which sense just the direction or just the time of arrival. It is also applicable to equipment which observes a Doppler shift between the transmitted and received signals. The term \"radar\" thus includes range finders, direction finders, and speed measuring equipment.\nThe invention particularly relates to radar systems employing pulse compression. In such a system a relatively long pulse is modulated before transmission, e.g., by sweeping its frequency from one value to another between the beginning and end of the pulse. On reception, the modulation of the pulse allows it to be compressed into a short time period. The use of pulse compression allows the transmission of a pulse which, because of its relatively long duration, has a relatively high energy whilst avoiding a loss of range resolution which would ordinarily be expected with an increase in pulse length.\nWhen a pulse is received after reflection from two features at slightly different ranges, the received echoes of this pulse may overlap in time. The pulse compression system is normally able to separate such overlapping echo signals. However, if the aforementioned features are at close range the received power may be so strong as to saturate initial stages of the receiver, causing it to behave non-linearly. This prevents the pulse compression system from operating properly and thus a close range target adjacent to another larger target or a large source of clutter may not be detected. This problem was discussed by Glasgow in U.S. Pat. No. 3,945,011 at the time when pulse compression was first being considered for use in civil aviation radars. The Glasgow Patent proposed a solution to the problem which involved the transmission of a short unswept pulse after each long swept pulse. The long pulses were used in the receiver after reflection from long range targets and were subjected to pulse compression whilst the short pulses were used when reflected from short range targets and were not subjected to pulse compression. The aforementioned problem was thus overcome.\nIt is a principal object of this invention to provide an alternative solution to the same problem as that dealt with by Glasgow and described above."}
{"text": "Various inorganic borates are known to be useful as fire retardants for materials such as wood, paper, cloth and organic polymers. For example, O'Shawnessy U.S. Pat. No. 3,897,387 describes the use of calcium or magnesium mineral borates as fire retardant agents for organic materials such as plastics, rubber and cellulosic material. Woods U.S. Pat. No. 3,816,307 described the use of Group I metal or ammonium borates having an average particle size of 0.1 and 25 microns as fire retardants for various plastic materials. Gouinlock U.S. Pat. No. 3,385,819 discloses the use of sodium tetraborate in combination with a polyhalogenated cyclopentadiene as a fire retardant for polymeric compositions. Woods et al U.S. Pat. No. 3,718,615 describes the use of a specific hydrated zinc borate as a fire retardant for halogen containing polymers. According to Schwarez et al. U.S. Pat. No. 3,560,441, halogenated polymers can be flame retarded by use of a glass formed by heating antimony trioxide with various inorganic materials including sodium borates. Japanese Patent Disclosure No. 1977-37947 describes preparation of a filler for fire resistant building materials by reacting sulfuric acid with borax in the presence of a porous inorganic refractory aggregate such as perlite. The resultant product is a refractory aggregate impregnated with a mixture of sodium sulfate and boric acid which can be used as a fire resistant filler for incorporating in polyurethane. Other patents such as Draganov U.S. Pat. No. 3,983,040, Blasius U.S. Pat. No. 4,156,654, Christianson et al. U.S. Pat. No. 4,172,804 and Quinto U.S. Pat. No. 4,173,666 describe the reaction of metal borates such as borax or the borate ores with sulfuric acid to form a granular composition for use as a fire retardant for cellulosic materials, especially insulation made from waste newspapers."}
{"text": "Finely powdered iron metal particles are pyrophoric as they can ignite and burn spontaneously upon direct contact with the air. The reaction is highly exothermic and self-sustainable, in which the iron particles react naturally with the oxygen in the atmosphere to form iron oxides, which is otherwise commonly known as the “rusting of steels”. The reactivity of the iron metal particles in the air has, however, been found to be largely a function of its particle sizes and the corresponding surface areas available to react with the oxygen in the atmosphere. Another important aspect of very fine iron particles, as with any nanosized metal powders, is its natural tendency to aggregate or sinter at elevated temperature to form larger particles resulting in dramatic loss of its pyrophoric properties.\nTherefore, significant efforts have been made over past decades to segregate and stabilize nanosized iron particles on various substrates with large surface areas, such as metallic foils, non-combustible fiber meshes, activated carbons, zeolites, and, most recently, carbon nanotubes, etc. The resulting pyrophoric materials, in various forms, have been an important subject in a wide array of applications such as cathode materials for fuel cells, active agent in chemical sensing device, and catalysts for ammonia synthesis, liquid-phase hydrogenation reactions for fine chemicals, and groundwater remediation, etc. Once passivated slightly in controlled environments, certain pyrophoric iron materials have found applications as pigments for magnetic tapes and in medical practice for its bacteriostatic properties. In military, pyrophoric iron materials have long been considered as the primary pyrotechnic charges in pyrophoric penetrator, ammunition training round markers, and infrared aerial decoy devices against heat-seeking missiles.\nPyrophoric iron particles found on various substrates are well documented, but most of those substrates reported in the literature do not actually provide the kind of confinement down to nanometer scales to prevent aggregation and sintering of fine particles at elevated temperature, as is the case with metallic foils and fabric materials. The pyrophoric iron composite materials thus obtained are normally dusty in nature, and the loading of active iron ingredients and the thermal activation are considered extremely delicate processes. In the cases of substrates with highly porous structures, such as zeolites or activated carbons, the transportation of the precursor iron ingredients into the existing narrow channels, less than 0.7 nm or around 1.0 nm of zeolites or activated carbons respectively, is severely hindered, so is its interaction with oxygen in air as pyrophoric materials. For instance, Gash et al, U.S. Patent Application Publication No. 2010/0139823, discloses methods for creating pyrophoric materials by first heating a carbon foam to create micropores, then depositing liquid solution containing metal ions into the micropores. This results in overall lower loading of pyrophoric iron particles and limited reactivity, and this class of materials is reported mostly as catalysts rather than pyrophoric materials.\nThe present invention provides an alternative approach to make pyrophoric foam materials by disclosing a simple in-situ or one-pot process for making pyrophoric foam materials in which the precursor chemicals with potential to form pyrophoric metal particles and porous carbon materials, respectively, are uniformly mixed, with or without a solvent, to produce a paste, molded or casted into any desirable geometry. Further thermal treatments results in the formation of the pyrophoric foams with pyrophoric metal particles uniformly distributed in a highly microporous carbon matrix which are also pyrophoric in nature."}
{"text": "An active matrix substrate is known that has the following configuration: a plurality of scanning lines (gate lines) extending in a first direction, and a plurality of data lines (source lines) extending in a second direction that is different from the first direction, are arranged on a substrate, and at the same time, transistors are arranged in correspondence to points of intersection between the scanning lines and the data lines (see Patent Document 1)."}
{"text": "1. Technical Field\nThe present invention relates generally to switch-reluctance motors (SRM's), and more specifically relates to a system and method of controlling the turn-on and turn-off angles of an SRM.\n2. Related Art\nElectric machines can be broadly classified into two categories on the basis of how they produce torque—electromagnetically or by variable reluctance. In the first category, motion is produced by the interaction of two magnetic fields, one generated by the stator and the other by the rotor. Two magnetic fields, mutually coupled, produce an electromagnetic torque tending to bring the fields into alignment. The same phenomenon causes opposite poles of bar magnets to attract and like poles to repel. The vast majority of motors in commercial use today operate on this principle.\nIn the second category, motion is produced as a result of the variable reluctance in the air gap between the rotor and the stator. When a stator winding is energized, producing a single magnetic field, reluctance torque is produced by the tendency of the rotor to move to its minimum reluctance position. As far as motors that operate on this principle, the switched-reluctance motor (SRM) falls into this class of machines.\nThe switched-reluctance motor (SRM) is under development for variable speed applications where the inherent characteristics of the SRM make commercial sense. To date, these applications include, e.g., sourcing aerospace power systems, starter/alternators for hybrid vehicles, and wind turbine applications. The aerospace and automotive applications are generally characterized by high-speed operation. The wind energy application is characterized by low speed, high torque operation.\nThe SRM produces torque through excitation that is synchronized to rotor position. The excitation is generally described by three excitation parameters: the turn-on angle θon, the turn-off angle θoff, and the reference current Iref. The turn-on angle θon refers to the rotor position where phase excitation (current) is initiated. The turn-off angle θoff refers to the rotor position where phase excitation is ended. The reference current Iref is the desired phase current during phase excitation. It should be noted that current continues to flow in the phase winding after θoff, but the phase current is monotonically decreasing by virtue of the state of the inverter switches. A computerized controller is typically utilized to control the parameters. Control of the excitation angles results in either positive net torque for motoring, or negative net torque for generating. A major challenge for implementing an SRM involves implementing a controller that effectively selects the turn-on and turn-off angles. Unfortunately, prior solutions have failed to provide a solution that allows for efficient operation of the SRM that is easily implemented. Such inefficiencies have lead to larger sizes, increased weight, and increased energy consumption.\nPrior art examples include U.S. Pat. No. 4,933,620, issued to MacMinn et al., which is hereby incorporated by reference and discloses a control strategy for low speed operation of an SRM. The invention deals with closed loop control of the turn-on angle of the SRM at low speeds and its hardware implementation. The suggested method tries to adjust the turn-on angle with the closed loop control such that first peak of the phase current occurs at the desired position. This control breaks down at the elevated speeds where the turn-on angle has no effect on the place of the first peak of the phase currents. The invention also keeps the conduction angle constant and advances the turn-off angle as the turn-on angle advances, but the conduction angle also needs to be adjusted as the operating conditions change for efficient system operation.\nU.S. Pat. No. 6,046,561 issued to Zup et al., which is hereby incorporated by reference, discloses a commutation control method for the SRM. The invention focuses on a hardware implementation of an approach commonly known as “phase advance” with a low cost position encoder. Phase advance means that the phase conduction interval is shifted relative to the angular interval over which the phase inductance is increasing. In U.S. Pat. No. 6,046,561, the phase advancement is made using only speed information and the conduction angle is kept constant in the invention. This assumption causes deterioration in the SRM efficiency. The most efficient turn-on and turn-off angles are nonlinear functions of the motor speed and motor phase currents and subject to change also with operating parameters.\nUnfortunately, none of the prior art provides efficient control of the excitation angles of the SRM at all operating points without extensive machine modeling and simulation. Accordingly, a need exists for a controller that can achieve these goals with a simple and efficient implementation."}
{"text": "1. Field of the Invention\nThe invention relates to a terminal structure, and more particularly to a terminal structure of an electrical connector having two horizontal contact surfaces opposite to each other.\n2. Description of the Related Art\nReferring to FIG. 1, a conventional terminal structure of an electrical connector includes an extension 11, a connection 12, and a contact 13, wherein the plate surfaces of the extension 11 and the connection 12 are in the same vertical plane. The contact 13 has two elastic arms 14 opposite to each other with a gap therebetween. The two elastic arms 14 have contact surfaces 15 close to each other. In addition, the two elastic arms 14 further have two twisted portions 16 to make the plate surfaces of the two elastic arms 14 opposite to each other.\nThe two elastic arms 14 of the conventional structure have twisted portions 16, respectively, to make the plate surfaces of the two elastic arms 14 opposite to each other so that wider plate surfaces may be utilized to form the contact surfaces 15. Instead of side surfaces 19 of a male terminal 17, the two contact surfaces 15 may contact top and bottom surfaces 18 of the inserted male terminal 17. Because the side surfaces 19 of the male terminal 17 are the cut-off surfaces that are rougher and the top and bottom surfaces are the smooth metal plate surfaces, it is preferred to contact the top and bottom surfaces in order to protect the gold-plated layers of the contact surfaces 15.\nAlthough the prior art structure has the above-mentioned advantage, the following drawbacks will be caused.\n1. The plate surfaces of the extension 11 and connection 12 are the vertical surfaces, a row of continuous terminals with gaps cannot be manufactured.\n2. Because the two elastic arms 14 are twisted to make the plate surfaces thereof opposite to each other, it is labor-consuming and difficult to control the precision of the twisted elastic arms. Thus, the two elastic arms may be slanted, and the gap between the contact surfaces is too small or too large to cause normal contacts with the inserted terminals.\nAs shown in FIG. 2, the two elastic arms are slanted, wherein the phantom line represents the hole 20 on the connector. In this case, the male terminal 17 can only be in point-contact with the contact surfaces 15 but not in normal surface-contact with the contact surfaces 15. As shown in FIG. 3, the gap between the contact surfaces of the two elastic arms is too large, and the male terminal 17 cannot be in contact with the contact surfaces 15. As shown in FIG. 4, the gap between the contact surfaces of the two elastic arms is too small, and the male terminal 17 cannot be easily inserted.\nReferring to FIG. 5, another conventional terminal structure of an electrical connector includes an extension 21, a connection 22, and a contact 23. The plate surface of the extension 21 is a horizontal surface and the plate surface of the connection 22 is a vertical surface. The contact 23 has two elastic arms 24 with a gap therebetween. The two elastic arms 24 have contact surfaces 25 close to each other. In addition, the two elastic arms 24 have twisted portions 26, respectively, to make the plate surfaces of the elastic arms 24 opposite to each other.\nAlthough the above-mentioned terminal structure may be utilized to form continuous terminals with gaps and facilitate the assembly, it also has the following drawbacks.\n1. Because the contact surfaces 25 of the two elastic arms 24 are vertical surfaces, they are in contact with the rough side surfaces 19 of the inserted male terminal 17, and the gold-plated layers of the contact surfaces 25 tends to be scratched.\n2. Because the two elastic arms 24 are respectively twisted to make the plate surfaces opposite to each other, it is labor-consuming and the precision of the two twisted elastic arms cannot be easily controlled."}
{"text": "Integrated circuits are made up of millions of active devices formed in or on a substrate, such as a silicon wafer. The active devices are chemically and physically connected into a substrate and are interconnected through the use of multilevel interconnects to form functional circuits. Typical multilevel interconnects comprise a first metal layer, an interlevel dielectric layer, and sometimes a third and subsequent metal layer. Interlevel dielectrics, such as doped and undoped silicon dioxide (SiO2) and/or low-κ dielectrics, are used to electrically isolate the different metal layers.\nThe electrical connections between different interconnection levels are made through the use of metal vias. U.S. Pat. No. 5,741,626, for example, describes a method for preparing dielectric TaN layers. Moreover, U.S. Pat. No. 4,789,648 describes a method for preparing multiple metallized layers and metallized vias in insulator films. In a similar manner, metal contacts are used to form electrical connections between interconnection levels and devices formed in a well. The metal vias and contacts may be filled with various metals and alloys, such as, for example, titanium (Ti), titanium nitride (TiN), aluminum copper (Al—Cu), aluminum silicon (Al—Si), copper (Cu), tungsten (W), platinum (Pt), ruthenium (Ru), iridium (Ir), and combinations thereof (hereinafter referred to as “via metals”).\nThe via metals generally employ an adhesion layer (i.e., a barrier film), such as a titanium (Ti), titanium nitride (TiN), tantalum (Ta), tantalum nitride (TaN), tungsten (W), or tungsten nitride (WN) barrier film, to adhere the via metals to the SiO2 substrate. At the contact level, the barrier film acts as a diffusion barrier to prevent the via metals from reacting with SiO2.\nIn one semiconductor manufacturing process, metal vias and/or contacts are formed by a blanket metal deposition followed by a chemical-mechanical polishing (CMP) step. In a typical process, via holes are etched through an interlevel dielectric (ILD) to interconnection lines or to a semiconductor substrate. Next, a barrier film is formed over the ILD and is directed into the etched via hole. Then, a via metal is blanket-deposited over the barrier film and into the via hole. Deposition is continued until the via hole is filled with the blanket-deposited metal. Finally, the excess metal is removed by a CMP process to form metal vias. Processes for the manufacturing and/or CMP of vias are disclosed in U.S. Pat. Nos. 4,671,851, 4,910,155, and 4,944,836.\nTypical metal CMP systems contain an abrasive material, such as silica or alumina, suspended in an oxidizing, aqueous medium. U.S. Pat. No. 5,244,534, for example, discloses a system containing alumina, hydrogen peroxide, and either potassium or ammonium hydroxide, which is useful in removing tungsten with little removal of the underlying insulating layer. U.S. Pat. No. 5,209,816 discloses a system useful for polishing aluminum that comprises perchloric acid, hydrogen peroxide, and a solid abrasive material in an aqueous medium. U.S. Pat. No. 5,340,370 discloses a tungsten polishing system comprising potassium ferricyanide, potassium acetate, acetic acid, and silica. U.S. Pat. Nos. 5,391,258 and 5,476,606 disclose systems for polishing a composite of metal and silica including an aqueous medium, abrasive particles, and an anion, which controls the rate of silica removal. U.S. Pat. No. 5,770,095 discloses polishing systems comprising an oxidizing agent, a chemical agent, and an etching agent selected from aminoacetic acid and amidosulfuric acid. U.S. Pat. No. 6,290,736 discloses a polishing composition for polishing a noble metal surface comprising an abrasive, a halogen compound, and an aqueous basic solution. Other polishing systems for use in CMP processes are described in U.S. Pat. Nos. 4,956,313, 5,137,544, 5,157,876, 5,354,490, and 5,527,423.\nBarrier films of titanium, titanium nitride, and like metals, such as tungsten, have a chemical activity similar to that of the via metals. Consequently, a single system can be used effectively to polish both Ti/TiN barrier films and via metals at similar rates. Ta and TaN barrier films, however, are significantly different from Ti, TiN, and like barrier films. Ta and TaN are relatively inert in chemical nature as compared to Ti and TiN. Accordingly, the aforementioned systems are significantly less effective at polishing tantalum layers than they are at polishing titanium layers (e.g., the tantalum removal rate is significantly lower than the titanium removal rate). While via metals and barrier metals are conventionally polished with a single system due to their similarly high removal rates, joint polishing of via metals and tantalum and similar materials using conventional polishing systems results in undesirable effects, such as oxide erosion and via metal dishing.\nSimilar problems with oxide erosion are observed when noble metals are used as the via metal. Noble metals have significantly lower chemical activity and are not adequately polished by conventional CMP compositions. Efficient planarization of noble metals often requires a CMP composition with an alkaline pH, resulting in undesirably higher removal rates of the oxide layer.\nConsequently, there remains a need for a system, composition, and/or method of polishing a substrate comprising a first metal layer and a second layer in a manner such that planarization efficiency, uniformity, and removal rate of the first metal layer are maximized and planarization of the second layer is minimized, thereby minimizing undesirable effects, such as first metal layer dishing, surface imperfections, and damage to underlying topography. The invention provides such a system, composition, and method. These and other characteristics and advantages of the invention will be apparent from the description of the invention provided herein."}
{"text": "The invention relates to a take-off apparatus for blown plastic sheets, produced by an extruder using a blow molding process with a collapsing device for the blown sheet, which is supplied with an expanded circular cross section, and with at least one rotatably mounted deflecting roller as well as with at least one non-rotatingly supported air tuning bar for transporting the collapsed blown sheet further.\nFor take-off apparatuses of this type, which can be constructed reversibly for distributing thickness tolerances, frozen into the blown sheet, over the whole width of the winder of the reeled film that is to be produced, the deflecting roller as well as the air turning bar usually are constructed cylindrically. During the operation of the take-off apparatus, compressed air emerges through the air outlet openings of the air turning bar, which allows the collapsed blown sheet, which is being brought along, to be suspended on a cushion of air as it loops around the air turning bar, in order to achieve frictionless guidance of the blown sheet about the non-rotating turning bar. If the pressure of the air emerging and, with that, of the cushion of air, is set too low, the blown sheet comes into contact with the air turning bar as it wraps around it, as a result of which scratches and similar damage appear on the blown sheet. In the case of blown sheets with strongly adhering surfaces, the latter may even adhere regionally, resulting in an interruption to the production process.\nFor this reason, efforts are made to work with the strongest possible cushion of air at the turning bar or bars, so that scratches on the surface of the sheet and interruptions to the production by adhesion of the blown sheet at an air turning bar are avoided. In this connection, however, it has been noted that the collapsed blown sheet, in the case of a strong cushion of air, inflates greatly on the reverse side of the air turning bar, that is, in the running direction of the sheet, behind the air nozzle region looped by the blown sheet after the latter has left the zenith of the air turning bar, with the result that the side folding edges of the collapsed blown sheet approach one another a little. Subsequently, if the blown sheet, which is narrower because the lateral folding edges are mutually approaching one another, runs onto a downstream, cylindrical deflecting roller, folding of the central region of the blown sheet occurs during the flat contact looping of the deflecting roller. This, in turn, has a negative effect on the quality of the blown sheet produced. In addition, due to the strong inflation of the blown sheet on the reverse side of the air turning bar, strong leakage losses of compressed air occur, which can be compensated for only with difficulty, since an increased ejection of compressed air in the looping region of the air turning bars can lead to unsteady running of the blown sheet with negative effects on the production process.\nFurthermore, take-off apparatuses have already become known, for which the deflecting roller and/or the air turning bar has a diameter, which differs over its length, for bringing about a concave or convex surface or peripheral surface of the deflecting roller and/or of the air turning bar. Admittedly, due to this shape of the peripheral surface, adhesion of the blown sheet in the looping region is counteracted by the cushion of air applied. However, a convex or concave arc-shaped cross-sectional profile is impressed on the blown sheet on the reverse side of a correspondingly constructed air turning bar. This, in turn, causes the side folded edges of the collapsed blown sheet to approach one another with the danger that folds will be formed in the middle region of the sheet during the further transport of the latter. In addition, there is a considerable leak of compressed air on the reverse side of the air turning bar in the middle or lateral regions of the latter, depending on whether this bar is equipped with a concave or convex peripheral surface.\nIt is an object of the invention to provide a take-off apparatus of the generic type, for which, while maintaining a sufficiently strong cushion of air to avoid adhesion of the collapsed blown sheet at the air turning bar in the looping region of the latter, high losses of compressed air due to leakage behind the air turning bar and the formation of an undesirable fold in the central region of the blown sheet during the further transport of the latter are avoided.\nPursuant to the invention, this objective is accomplished by providing a configuration of a take-off apparatus for plastic blown sheets, produced by an extruder in the blow molding process, with a collapsing device for the blown sheet, which is supplied with an expanded circular cross section, and with at least one rotatably mounted deflecting roller as well as with at least one non-rotatingly supported air turning bar for transporting the collapsed blown sheet further, the air turning bar, in a region of its peripheral surface, looped by the collapsed blown sheet during its transport, being provided with an arrangement of air outlet openings, characterized in that the air turning bar, with a constant external diameter, has a course, which is arc-shaped in the axial direction and defines a plane of curvature with a concave side and a convex side of the air turning bar, and in that the air outlet openings are disposed on the concave side of the air turning bar. Owing to the fact that, in the case of this configuration, the air outlet openings of the air turning bar are disposed on the concave side of the latter, that is, on the side facing the chord of an imaginary, associated circular section, the space is created, which is required for applying a sufficiently large cushion of air, which avoids adhesion of the collapsed blown sheet, as the latter loops around the air turning bar, even when the surface has an adhesive nature. At the same time, the external diameter of the air turning bar, which is constant throughout, ensures that the blown sheet, at the zenith of the air turning bar, runs off from this in a straight plane over the convex side of the air turning bar, so that an inflation of the blown sheet on the reverse side of the air turning bar by entrained compressed air from the looping region, in conjunction with an inwards deformation of the peripheral surface of the turning bar, is counteracted. With that, the cause for an undesirable fold formation in the central region of the blown sheet, during the further transport of the latter, is also excluded, so that the prerequisites for a defect-free production of blown film, even of those with strongly adhering surfaces, are created.\nIn general, the plane of curvature of the air turning bar is disposed parallel to a transporting plane of the collapsed blown sheet, which is specified by a deflecting roller disposed upstream. In the case of a take-off apparatus, which is commonly used and which takes up the collapsed sheet, transported perpendicularly upwards from the extruder, this means a horizontal arrangement of the deflecting roller or rollers and, correspondingly, a horizontal arrangement of the air turning bar or bars for transporting the collapsed blown sheet by the take-off apparatus over conveying path sections, which in turn are horizontal.\nAn adjustable support of the air tuning bar or bars advantageously is such that its plane of curvature can be adjusted from a basic position, which is parallel to the transporting plane of the collapsed blown sheet and usually horizontal, into a position inclined to this at an angle of about 10xc2x0 to 20xc2x0. By these means, the cushion of air can be regulated and the course of the blown sheet optimized.\nThe height of the arc, presented by the air turning bar, that is, the measure of the distance between the apex of the arc and the chord of an imaginary, associated circular section, depends on the length of the air turning bar, which in turn depends on the maximum width of the blown sheet produced. The arc height is only a few millimeters and usually ranges from about 1 to 10 mm.\nIn order to avoid losses of compressed air and factors that interfere with the transporting behavior of the blown sheet in the looping region of the air-guiding roller when transporting narrow blown sheets, the width of which is less than the length of turning bars designed for a blown sheet of maximum width, the air outlet openings, which are not required in the looping region of the air turning bar outside of the transporting region, can be closed off, for example, by a diaphragm device, mounted in the interior of the air turning bar.\nFurther details and advantageous effects of the invention arise from the description of the object of the invention below, in conjunction with the drawing, in which an example of a take-off apparatus of the present invention is illustrated diagrammatically."}
{"text": "1. Field of the Invention\nThe present invention generally relates to racking systems or dishracks for dishwashing appliances and in particular to dishwasher racks having pivotable fence assemblies which enhance the loading flexibility of the racks.\n2. Description of the Related Art\nIn the design of dishwasher racks it has generally been industry practice to design racks with dish support members arranged for loading xe2x80x9cstandardxe2x80x9d dish sizes. It is, however, likely that most households will include some xe2x80x9cnon-standardxe2x80x9d dishes as well as various odd shaped cooking utensils which will not load properly in a dishwasher rack designed for only xe2x80x9cstandardxe2x80x9d dishes. To alleviate the problem presented by xe2x80x9cnon-standardxe2x80x9d dishes and other cooking utensils, several rack designs have evolved which have included movable or adjustable portions.\nGeiger et al, in U.S. Pat. No. 3,269,548 issued Aug. 30, 1966, teach a rack system which is vertically adjustable within the dishwasher to accommodate large articles. The upper rack also includes several pivoted divider members. Each of the pivoted divider members is formed of wire members and is supported on the upper rack by a pair of hinge clips secured to the base of the wire rack for supporting the ends of the pivoted divider members. A separate sheet metal detent member is positioned at the end of the rack for holding the pivoted divider members in the upright position.\nSmith, in U.S. Pat. No. 3,402,975 issued Sep. 24, 1968, teaches a top loading dishwasher having two levels of dishracks. Access to the lower dishrack is attained by utilizing a horizontally movable rack in cooperation with a hinged shelf which is pivotally supported between a pair of the dishwasher side walls at one end and by a portion of the movable rack at the other end. The hinged shelf includes a pair of pivotal divider assemblies each of which may be maintained in an upright posture by a detent assembly at one end of the divider assembly.\nFiocca et al, in U.S. Pat. No. 3,752,322 issued Aug. 14, 1973, teach a dishwasher rack construction in which a pivoted element can be retained in one of a plurality of intermediate positions by an arcuate toothed member located at the front of the rack.\nYake, in U.S. Pat. No. 4,046,261 issued Sep. 6, 1977, discloses a dishwasher rack which includes a pair of brackets fixed to the bottom of the rack and having molded slots and retainers for receiving a plurality of fences in a plurality of positions. These fences are either vertically mounted between the pair of brackets or are removed from the dishwasher rack. There is no pivotal movement of individual fences but only selective use of these fences for adjusting the distance between adjacent fences and between the rack side walls.\nPrior patents show that considerable effort has been focused on the development of dishwasher racks and on the development of dishwasher racks which include flexible or variable spacing arrangements for accommodating a variety of dish loads.\nIt is therefore an object of the instant invention to provide an improved dishwasher rack with pivotable fences.\nIt is a further object of the instant invention to provide a dishwasher rack with a plurality of fences which are pivoted about a common axis.\nIt is a still further object of the instant invention to provide a dishwasher rack with a pivotable fence wherein the pivotable fence is positively secured to the dishwasher rack in a unique and simple manner.\nAccording to the present invention, the foregoing and other objects are attained a dishwasher rack which includes a basket having a plurality of wires forming a bottom wall, a pair of side walls and a front and back wall. The bottom wall has a plurality of upwardly oriented tines connected to and extending from the bottom wall in a first area and further has an open area without upwardly oriented tines. A first fence is provided having an axle wire and a plurality of support wires. The axle wire is rotatably supported adjacent the bottom wall such that the first fence may be rotated between a front position and a back position. A second fence is provided having a connection wire and a plurality of support wires. The connection wire is rotatably connected to axle wire such that the second fence may rotate about the axle wire between a front position and a back position. The first fence and second fence may be rotated to the first position such that the support wires of the first fence and the support wires of the second fence form an upstanding array for supporting dishes. Alternatively, one or both of the first fence and second fence may be rotated to the back position.\nThe first fence is connected to the basket without the need for additional clips. In particular, the ends of the axle wire are captured in sockets formed by the combination of the basket and wheel assemblies. The second fence may connect to the first fence without the use of clips or, alternatively, with the use of clips."}
{"text": "1. Field of the Invention\nThe present invention generally relates to battery cover assemblies and, more particularly, to a battery cover assembly for use in a portable electronic device.\n2. Discussion of the Related Art\nBatteries are widely used in portable electronic devices such as personal digital assistants (PDAs) and telephones. A battery is usually removably mounted in a receptacle of a housing of an electronic device. The receptacle is generally covered by a battery cover. When the battery is damaged, dead, or in need of recharging, the cover is removed from the housing, whereupon the battery can be taken out and later replaced.\nA fixing/attachment mechanism is generally employed in a conventional battery cover, to engage with a housing of the portable electronic device. For example, a battery cover of an Alcatel® OT310 telephone includes a fixing mechanism. The fixing mechanism incorporates a pair of hooks arranged at a lower end of the battery cover and a locking pin arranged at an upper end thereof. Correspondingly, a pair of slots is defined at a lower end of a backside of the housing, and a locking hole is defined at an upper end of the back side. In assembly, the hooks are inserted into the corresponding slots. Then, the battery cover is pressed downwardly such that the locking pins of the battery cover are inserted into the corresponding locking holes of the housing. The battery cover is thus assembled to the housing of the telephone. The battery cover is simple in structure, and the engagement between the battery cover and the housing of the telephone is secure. However, during disassembly of the cover from the housing of the telephone, the battery cover is liable to be damaged because great force has to be exerted on the battery cover to detach the battery cover from the housing. As a result, it can be inconvenient for a user to change a battery.\nWhat is needed, therefore, is a new battery cover assembly for a portable electronic device which makes it convenient for a user to detach a battery cover from a housing of the portable electronic device."}
{"text": "This invention generally relates to medical devices, and particularly to intracorporeal devices for therapeutic or diagnostic uses, such as balloon catheters. In percutaneous transluminal coronary angioplasty (PTCA) procedures, a guiding catheter is advanced until the distal tip of the guiding catheter is seated in the ostium of a desired coronary artery. A guidewire, positioned within an inner lumen of a dilatation catheter, is first advanced out of the distal end of the guiding catheter into the patient's coronary artery until the distal end of the guidewire crosses a lesion to be dilated. Then the dilatation catheter having an inflatable balloon on the distal portion thereof is advanced into the patient's coronary anatomy, over the previously introduced guidewire, until the balloon of the dilatation catheter is properly positioned across the lesion. Once properly positioned, the dilatation balloon is inflated with fluid one or more times to a predetermined size at relatively high pressures (e.g. greater than 8 atmospheres) so that the stenosis is compressed against the arterial wall and the wall expanded to open up the passageway. Generally, the inflated diameter of the balloon is approximately the same diameter as the native diameter of the body lumen being dilated so as to complete the dilatation but not overexpand the artery wall. Substantial, uncontrolled expansion of the balloon against the vessel wall can cause trauma to the vessel wall. After the balloon is finally deflated, blood flow resumes through the dilated artery and the dilatation catheter can be removed therefrom.\nIn such angioplasty procedures, there may be restenosis of the artery, i.e. reformation of the arterial blockage, which necessitates either another angioplasty procedure, or some other method of repairing or strengthening the dilated area. To reduce the restenosis rate and to strengthen the dilated area, physicians frequently implant a stent inside the artery at the site of the lesion. Stents may also be used to repair vessels having an intimal flap or dissection or to generally strengthen a weakened section of a vessel. Stents are usually delivered to a desired location within a coronary artery in a contracted condition on a balloon of a catheter which is similar in many respects to a balloon angioplasty catheter, and expanded to a larger diameter by expansion of the balloon. The balloon is deflated to remove the catheter and the stent left in place within the artery at the site of the dilated lesion. Stent covers on an inner or an outer surface of the stent have been used in, for example, the treatment of pseudo-aneurysms and perforated arteries, and to prevent prolapse of plaque. Similarly, vascular grafts comprising cylindrical tubes made from tissue or synthetic materials such as polyester, expanded polytetrafluoroethylene, and DACRON may be implanted in vessels to strengthen or repair the vessel, or used in an anastomosis procedure to connect vessels segments together.\nIn the design of catheter balloons, characteristics such as strength, compliance, and profile of the balloon are carefully tailored depending on the desired use of the balloon catheter, and the balloon material and manufacturing procedure are chosen to provide the desired balloon characteristics. A variety of polymeric materials are conventionally used in catheter balloons. Use of polymeric materials such as PET that do not stretch appreciably consequently necessitates that the balloon is formed by blow molding, and the deflated blow molded balloon forms wings which are folded around the catheter shaft prior to inflation of the balloon in the patient's body lumen. However, it can be desirable to employ balloons, referred to as formed-in-place balloons, that are not folded prior to inflation, but which are instead expanded to the working diameter within the patient's body lumen from a generally cylindrical or tubular shape (i.e., essentially no wings) that conforms to the catheter shaft.\nA catheter balloon formed of expanded ultra high molecular weight polyethylene (eUHMWPE) has been suggested. However, eUHMWPE has proven difficult to bond to balloon liner materials and to catheter shaft materials.\nIt would be a significant advance to provide a catheter balloon, or other medical device component, with improved performance and bondability."}
{"text": "This invention relates to an organopolysiloxane composition for preparing peelable film. More specifically, this invention relates to a solventless type organopolysiloxane composition for preparing peelable film which exhibits moderate peelability from various adhesive materials whose properties do not depend on time and whose long-lasting adhesive capacity is high.\nIt is an established fact that a cured film of organopolysiloxane on the surface of a flexible material such as paper, laminated paper, synthetic film, knitted goods or metal foil can increase the peelability of the flexible material from an adhesive. Examples of the organopolysiloxane composition used for the above-mentioned purpose are polycondensation type organopolysiloxane compositions composed of a diorganopolysiloxane having terminal hydroxyl groups, an organohydrogenpolysiloxane and organotin compounds as curing catalysts (Japanese Pat. No. Sho 35 {1960}-13709); and addition-reaction type organopolysiloxane compositions composed of a vinyl-containing organopolysiloxane, an organohydrogenpolysiloxane and platinum compounds as curing catalysts (Japanese Pat. No. Sho 46 {1971}-26798). Japanese Pat. No. Sho 49 {1974}-27033 and Japanese Kokai Pat. No. Sho 52 {1977}-86985 describe an organopolysiloxane composition for preparing peelable film which exhibits a moderate peelability. However, these organopolysiloxane compositions are principally composed of organopolysiloxanes with a high degree of polymerization so that they must be diluted with a large quantity of organic solvent in order to be coated on the base material effectively. The organic solvents include aliphatic solvents such as hexane and heptane, aromatic solvents such as benzene, toluene and xylene, halogenated hydrocarbon solvents such as trichloroethylene, perchloroethylene and carbon tetrachloride, and ethyl acetate. However, these solvents are toxic to workers, cause environmental pollution, fires and explosions, and require recovery equipment. These are the traditional problems to be solved.\nIn order to eliminate the defects of the traditional method, Japanese Kokai Pat. No. Sho 47 {1972}-32072 and Japanese Pat. No. Sho 53 {1978}-39791 describe an organopolysiloxane for preparing peelable film which can be cured by an addition reaction, which does not use an organic solvent, that is, it can be coated without diluting it with solvent. These solventless type organopolysiloxane compositions form easily peelable films so that they cannot be used for processing paper and wrapping paper for asphalt which require a moderate peeling capacity and for such materials as tapes and labels which also require a moderate peeling capacity."}
{"text": "This invention relates to a stencil sheet discharging device of a stencil printing machine, and more particularly, a stencil sheet discharging device of a stencil printing machine in which the discharged stencil sheets fed into a discharged stencil sheet box are compressed by a discharged stencil sheet compression plate and filled in it.\nAs shown in FIG. 8A, in this type of stencil sheet discharging device in the prior art, a stencil sheet to be discharged is fed into a stencil sheet discharging box 2 by a discharged stencil sheet feeding operation under a state in which a discharging stencil sheet compression plate 1 is positioned at a home position. Upon completion of the discharged stencil sheet feeding operation, the discharged stencil sheet compression plate 1 descends as shown in FIG. 8B to start a compressing operation. At this time, a compression sensor (a photo-interrupter) 3 composed of a light emitting section and a light receiving section is changed from its light receiving state to a light shielding state by a compression sensor plate 5 arranged at a driving shaft 4. In addition, a torque is added to the discharged stencil sheet compression plate 1 from the driving shaft 4 through a coil spring (not shown).\nThen, when the discharged stencil sheet compression plate 1 further descends and the compression sensor 3 is changed from the light shielding state to the light receiving state as shown in FIG. 8C, the compression motor is turned off. Then, after elapsing a specified period of time, the discharged stencil sheet compression plate 1 is lifted. At this time, as shown in FIG. 8C, fill-up sensor (a photo-interrupter) 6 comprised of a light emitting section and a light receiving section is shielded against light by a fill-up sensor plate 7.\nIn turn, as shown in FIG. 8D, when a stencil sheet discharged into the discharged stencil sheet box 2 fills up, the discharged stencil sheet compression plate 1 cannot descend by the discharged stencil sheet. Due to this fact, only the driving shaft 4 (the compression sensor plate 5) is rotated. Accordingly, even if the compression sensor 3 becomes a light receiving state again after it became a light shielding state, the fill-up sensor 6 is not shielded against light by a fill-up sensor plate 7. With such an arrangement as above, it is detected that the discharged stencil sheet box 2 is full of discharged stencil sheets.\nThe aforesaid stencil sheet discharging device for the stencil printing machine in the prior art can detect by the fill-up sensor 6 that the discharged stencil sheet box 2 is full of stencil sheets. However, this cannot detect the residual amount (vacant capacity) in the discharged stencil sheet box 2. Thus, this device shows a problem that the time required for filling the box with stencil sheets cannot be estimated.\nIn addition, there occurs occasionally that an operator desires to adjust the fill-up position in the discharged stencil sheet box 2 in reference to its state of use. However, a fixing position of the fill-up sensor 6 in the prior art stencil sheet discharging device for the stencil printing machine must be adjusted. Accordingly, this type of device has a problem that a fixing position of the fill-up sensor 6 cannot be easily adjusted.\nFurther, in the case that much amount of residual discharged stencil sheets is present in the box 2, all of the discharged stencil sheets is not compressed. Then, as the residual amount of discharged stencil sheet box 2 is reduced, a compression force is gradually applied to each of the discharged stencil sheets. Accordingly, it shows some problems that an appropriate compression force cannot be applied for each of the discharged stencil sheets and a superior compression cannot be carried out.\nThe present invention has been invented to solve the problems found in the prior art and it is an object of the present invention to reduce the number of sensors for use in sensing the compressing position or the compression force of the discharged stencil sheet compression plate and to reduce cost.\nIt is another object of the present invention to estimate a time when the box is full of stencil sheets by detecting the residual amount in the discharged stencil sheet box.\nIt is a still further object of the present invention to adjust a fill-up sensor fixing position in an easy manner.\nIt is a still yet further object of the present invention to apply a proper compression force for every discharged stencil sheet in the discharged stencil sheet box and to perform a superior compression.\nIn order to accomplish the aforesaid objects, according to a first aspect of the present invention, there is provided a stencil sheet discharging device for a stencil printing machine comprising:\na discharged stencil sheet compression plate 43 for use in compressing a discharged stencil sheet fed into a discharged stencil sheet box 42 within the discharged stencil sheet box 42;\na driving means (a compression motor 44) for driving the discharged stencil sheet compression plate 43 and compressing the discharged stencil sheet within the discharged stencil sheet box 42;\na compression force detecting means for sensing that a predetermined compression force is added to the discharged stencil sheet;\na control means for stopping a compressing operation performed by the driving means 44 when a predetermined compressive force is detected by said compression force detecting means 50;\na moving amount detecting means for detecting a moving amount of a discharged stencil sheet compression plate 43 when a predetermined compressive force is detected by said compression force detecting means; and\na display means (a discharged stencil sheet amount display unit 94) for displaying a residual amount in the discharged stencil sheet box 42 in reference to a moving amount detected by the moving amount detecting means.\nWith this stencil sheet discharging device, it becomes possible to estimate a time in which the box is filled with the discharged stencil sheets due to the fact that a moving amount of the discharged stencil sheet compression plate 43 is detected and a residual amount in the discharged stencil sheet box 42 is displayed. In addition, it becomes possible to perform a superior compression and store much amount of discharged stencil sheets in the discharged stencil sheet box 42 due to the fact that every discharged stencil sheet is compressed with a predetermined compression force.\nAccording to a second aspect of the present invention, there is provided a stencil sheet discharging device for a stencil printing machine comprising:\na discharged stencil sheet compression plate 43 for use in compressing a discharged stencil sheet fed into a discharged stencil sheet box 42 within the discharged stencil sheet box 42;\na driving means (a compression motor 44) for driving the discharged stencil sheet compression plate 43 and compressing the discharged stencil sheet within the discharged stencil sheet box 42;\nan encoder 50 for generating the number of pulses corresponding to the moving amount of the discharged stencil sheet compression plate 43;\na compression force detecting means for detecting a pulse width of pulse generated from the encoder 50 and detecting that a predetermined compression force is added to the discharged stencil sheet when said pulse width exceeds a predetermined value;\na control means for stopping a compressing operation performed by the driving means 44 when a predetermined compression force is detected by the compression force detecting means;\na moving amount detecting means for counting pulses generated from an encoder 50 until a predetermined compression force is detected by the compression force detecting means and detecting a moving amount of the discharged stencil sheet compression plate in response to the counted number of pulses; and\na display means (a discharged stencil sheet amount display unit 94) for displaying a residual amount in the discharged stencil sheet box 42 in reference to a moving amount detected by the moving amount detecting means.\nThe invention according to a second aspect of the present invention utilizes one unit of encoder as a sensor for detecting either a compressing position or a compression force of the discharged stencil sheet compression plate as compared with the invention according to a first aspect of the present invention so as to reduce cost of the sensor.\nAs indicated in a third aspect of the present invention, there is provided a home position sensor 56 for use in sensing a home position of the discharged stencil sheet compression plate 43, the moving amount sensing means either resets or pre-sets the number of pulses to be counted when the home position sensor 56 detects the home position of the discharged stencil sheet compression plate.\nFurther, as indicated in a fourth aspect of the present invention, there is provided a setting means for setting a compression fill-up position of the discharged stencil sheet box 42, the display means 94 displays the residual amount in the discharged stencil sheet box 42 in response to the compression fill-up position set by the setting means and the moving amount detected by the moving amount detecting means. That is, it is possible to set the compression fill-up position in the discharged stencil sheet box 42 in an easy manner and further the display means 94 can display the residual amount in response to the set fill-up position.\nAs indicated in a fifth aspect of the present invention, said control means controls said driving means 44 so that said discharged stencil sheet compression plate 43 is operated to compress the discharged stencil sheet within said discharged stencil sheet box 42 until the predetermined compression force is detected by said compression force detecting means, said discharged stencil sheet compression plate 43 is stopped at the compressing position for a predetermined period of time in case where the predetermined compression force is detected by said compression force detecting means, and said discharged stencil sheet compression plate 43 is returned back to the predetermined waiting position, i.e. home position."}
{"text": "Computer networks have become widely used in society today. Mobile nodes, for example, cellular telephones, are often used in conjunction with these networks. User devices, such as computers, may be connected to mobile nodes. In this way, the user device (e.g., the computer) may be in direct communication with some or all of the networks and other devices connected to these networks.\nAs mobile nodes move between networks, the ability to locate and communicate with the mobile node is preferably maintained. For instance, the mobile node may be assigned a home network. The home network may include a home agent. The mobile node may move to other networks (“foreign networks”) that are not the home network. As the mobile node moves to these foreign networks it may register a care-of-address with the home network and the home agent in the home network. The mobile node also may register this care-of-address with a foreign agent on the foreign network. In this way, messages can be transmitted from a device in any network to the mobile node, no matter where the mobile node is located. Also, the mobile node (and user devices coupled to the mobile node) may transmit messages to the home network and receive messages from the home network.\nA tunnel may be used to transmit the information from the tunnel endpoints. In tunneling, packets or frames from one network are placed inside frames (“encapsulated”) of another network. The encapsulated frames may include a header with sufficient routing information to transmit the encapsulated frame from a source to a destination. In one example of tunneling, the home agent may tunnel data to the mobile node when the mobile node is on a foreign network. In another example of tunneling, a tunnel may be established between the foreign agent and the home agent to transmit data between the mobile node and the home agent when the mobile node is attached to a foreign network. Once the encapsulated frames reach the tunnel endpoint, they are unencapulated and transmitted to the ultimate destination."}
{"text": "The present invention relates to an analytical method for assessing the presence of alkaline sphingomyelinase in the stools or biological fluids of patients in need of such an assessment. The invention also relates to a kit for carrying out the analytical method.\nMore particularly the method of the present invention is an in vitro fluorometric method for detecting alkaline sphingomyelinase which, as will be described in detail hereinbelow, is a marker of serious pathological states such as colon cancer and familial adenomatous polyposis.\nThe enzyme sphingomyelinase (sphingomyelin phosphodiesterase, SMase) catalyzes the hydrolysis of sphingomyelin to ceramide and choline phosphate.\nThree different types of SMase (acidic, neutral and alkaline) have been identified to-date, which occur as several isoforms, as follows:\nlysosomal acidic SMase (A-SMase);\ncytosolic Zn2+-dependent acidic SMase;\nmembrane neutral magnesium-dependent SMase (N—SMase);\ncytosolic magnesium-independent N—SMase; and\nalkaline SMase.\nSMases have been shown to play a role in a wide variety of physiologic and pathological processes, including: lysosomal hydrolysis of endocytosed SM, ceramide mediated cell signaling, atherogenesis, terminal differentiation, cell cycles arrest, apoptosis, inflammation, and the regulation of eukaryotic stress responses.\nIn contrast to acidic and neutral SMase, which are currently present in cells as lysosomal and membrane-bound enzymes, respectively, alkaline SMase exhibits tissue and species difference. In human beings, the alkaline SMase is found in intestinal mucosa and bile. Alkaline SMase starts to appear in the duodenum, reaches a high level in the intestine, especially in the distal part of the jejunum, and occurs in considerable amounts in the colon and rectum. This SMase presents optimal alkaline pH at 9.0, is Mg2+-independent, bile salt-dependent and trypsin-resistant.\nThe pathological importance of alkaline SMase has only recently been recognized and this has prompted several studies to be carried out, mainly for the following reasons.\nFirst, the enzyme may be responsible for the hydrolysis of the dietary sphingomyelin occurring substantially in milk, eggs, meat and fish. Second, this enzyme may regulate cholesterol absorption. Third, the presence of alkaline SMase along the intestinal tract and its selective decrease detected in colorectal carcinoma suggests that this enzyme plays a role in intestinal carcinogenesis, since under physiological conditions, it stimulates apoptosis and protects the intestinal mucosa against carcinogenesis.\nPrevious studies have also shown that, under physiological conditions, alkaline SMase is dissociated by bile salts from intestinal mucosal membrane to the lumen. However, under pathological conditions, whereby bile salt concentration is abnormally increased, the dissociation of alkaline SMase by bile salts may exceed the biosynthesis of the enzyme, resulting in a low level of activity of alkaline SMase in the mucosa, and an abnormally increased excretion of the enzyme in the feces or in biological fluids, i.e., bile. Consequently, the excess of alkaline SMase excreted in the stools or in biological fluids over normal, basal values, may be interpreted as a valuable diagnostic marker for colon rectal carcinoma and familial adenomatous polyposis, hence; the need of a reliable assay for detecting alkaline SMase in the stools or in biological fluids of patients likely to be suffering from the aforesaid pathologies of the intestinal tract.\nIn addition, some bacteria strains (e.g., Streptococcus htermophilus or Lactobacilli) contain high levels of SMase, and the assessment of alkaline SMase may provide a method to evaluate changes in the number of said bacteria, i.e., after a treatment with probiotics or/and probiotic-based products.\nPrevious methods for assaying alkaline SMase are already known. The activity of the SMases can be determined either in vivo through cell labeled with a radioactive precursor of SM and then determining the labeling product levels or in vitro using radiolabeled SM or a chromogenic analog of SM or colored and fluorescent derivatives of neutral SM.\nThese known commonly used assays are not entirely satisfactory since they are potentially very hazardous insofar as they are radioactive assays and less sensitive than a fluorometric assay."}
{"text": "The present invention relates in general to voice mail systems and more specifically to a method and apparatus for identifying and extracting voiced telephone numbers and email addresses from recorded voice mail messages.\nA common complaint among users of voice mail systems relates to difficulties in extracting and transcribing important information from voice mail messages for later use. As an example, it is common for a caller whose call is answered by voice mail to leave a voice mail message that includes a telephone number or email address at which the caller may be contacted. Unless the caller enunciates the telephone number or email address very slowly, the voice mail recipient may be forced to replay all or a portion of the voice mail message in order to transcribe the telephone number or email address for later use. Most callers find very slow enunciation to be awkward and unnatural, with the result that the majority of such voice mail messages must be replayed once or even several times in order for the voice mail recipient to transcribe the information. Furthermore, the process for replaying voice mail messages is often complex (e.g. requiring memorization of keys on a telephone keypad which have been assigned for the replay function). The process can also be tedious (e.g. some systems require that the entire message be replayed in order to extract only the portion relating to the telephone number or email address).\nAccording to the present invention, a system is provided for automatically identifying enunciation of telephone numbers and email addresses in voice mail messages for presentation to a user via either a voice or text, without the user being required to replay the message. Furthermore, according to the invention the extracted telephone number or email addresses may be forwarded to an application for further processing (e.g. to dial the extracted number, add the email address and/or telephone number to a contacts database, etc.)\nAccording to another aspect of the invention, a user may specify text strings to be searched (other than telephone numbers and email addresses). Thus, a user can specify certain words such as xe2x80x9cextensionxe2x80x9d, xe2x80x9cemailxe2x80x9d, etc., in order to qualify the extracted number or email address information, or words such as xe2x80x9curgentlyxe2x80x9d or xe2x80x9cimmediatelyxe2x80x9d to identify important voice mail messages."}
{"text": "The present disclosure relates generally to medical connectors used in fluid transfer applications. More particular, it relates to a vial adapter for the transfer of fluids in medical settings without exposure of the fluid to an ambient atmosphere.\nMedical connectors are widely used to transmit, prepare, and deliver medical fluids. The preparation of a medical fluid may include the delivery, dilution, reconstitution, and withdrawal of a medical fluid or a component thereof with a container such as a vial.\nIn some instances, such as with chemotherapy treatment, the medical fluid is hazardous. Particularly, repeated exposure to the medical fluid, such as by medical personnel, is hazardous. An example instance of medical fluid transfer is the reconstitution of a medication. Reconstitution is often conducted within a sealed vial containing a medical fluid, or a constituent thereof, in any state of matter. This process requires a diluent to be delivered into the vial. However, delivery of the diluent into the sealed vial causes displacement of gas within the vial. If the gas were permitted to enter the ambient atmosphere, people within the ambient atmosphere may be exposed to the gas. In some instances, the medical fluid itself may be transmitted into the ambient atmosphere during reconstitution."}
{"text": "This application is based on application Ser. No. 2000-45858 filed in Japan on Feb. 23, 2000, the contents of which are hereby incorporated by reference.\nThe present invention relates to a process cartridge and an image forming apparatus using the same.\nIn an image forming apparatus used in a copying machine, a printer or the like, for the purposes of exchanging or repairing elements necessary for image forming processes, treating a jam or the like, process cartridges including the elements are detachably mounted on the image forming apparatus. Such process cartridges are necessary to be securely mounted and positioned on the image forming apparatus so as not to cause a maloperation.\nU.S. Pat. No. 5,066,976 proposes such a mechanism that if a process cartridge is not perfectly push into an apparatus body, a part of holding member for holding the process cartridge protrudes out of an opening portion of grip, whereby it can be possible to judge whether or not the process cartridge is pushed into a predetermined position. Japanese Utility Model Laid-open publication No. 63-100736 proposes such a mechanism that when a cover is closed, a stopper mechanism provided on the cover pushes and fixes a process unit. Japanese Utility Model Laid-open publication No. 62-149046 proposes such a mechanism that a process unit detachably mounted so as to approach to and separate from a photosensitive is moved by a lever to position in place.\nHowever, the positioning mechanisms of process cartridges proposed above have a disadvantage that when the process cartridges are positioned in a half-locked state or a half-inserted condition, or when the operation of the lever is omitted, the positioning of the process cartridges has to be conducted again, and also have a disadvantage that the push-fixation mechanism is complicated.\nAccordingly, it is an object of the present invention to provide a process cartridge which has a simple mechanism and can be easily positioned with respect to an apparatus body and securely mounted on the apparatus body even in a half-locked state or a half-inserted condition, and also provide an image forming apparatus using the process cartridge.\nAccording to a first aspect of the present invention, there is provided a process cartridge which can be mounted on and dismounted from an image forming apparatus, the process cartridge comprising:\na unit comprising an element used for an image forming process;\na positioning member having a positioning member which engages with an engagement portion formed on the image forming apparatus to position the process cartridge at a predetermined mount position in the image forming apparatus;\na holding member for holding the unit to allow displacement of the unit relative to the positioning member; and\na restricting and releasing mechanism for restricting or releasing the displacement of the unit relative to the positioning member.\nIn the first aspect of the present invention, when the process cartridge is mounted on the image forming apparatus, the unit is secured by restricting the displacement of the unit relative to the positioning member, it is possible to stably surely mount the process unit. The term xe2x80x9cdisplacementxe2x80x9d in this specification of the present invention includes pivot, vibration, movement, swing and so on.\nIt is preferable that the engagement portion of the image forming apparatus comprises a recess or projection, and wherein the positioning member comprises a projection or recess which engages with the recess or projection of the engagement portion. In this case, the projection or recess of the positioning member may comprise a plurality of projections or recesses which are disposed around the holding member.\nPreferably, the holding member comprises a screw member.\nPreferably, the process cartridge further comprises a cover member for interlocking with the restricting and releasing mechanism to cover or open a specific portion of the unit. In this case, the element may comprise a toner container and a developing roller, and cover member may cover or open the developing roller.\nPreferably, the restricting and releasing mechanism comes into contact with a member provided on the image forming apparatus and operates. In this case, the restricting and releasing mechanism may operate when the process cartridge is mounted on and dismounted from the image forming apparatus. Alternately, the restricting and releasing mechanism mat operate when the process cartridge is mounted on the image forming apparatus.\nAccording to a second aspect of the present invention, there is provided a process cartridge which can be mounted on and dismounted from an image forming apparatus, the process cartridge comprising:\na unit comprising an element used for an image forming process;\na positioning member having a first positioning member which engages with an engagement portion formed on the image forming apparatus to position the process cartridge at a first predetermined position in the image forming apparatus; and\na first pressed portion which is pressed by a first press portion provided on the image forming apparatus when the process cartridge is moved from the first predetermined position to a second predetermined position to position the process cartridge at the second predetermined position.\nIn the second aspect of the present invention, when the process cartridge moves in the image forming apparatus, the first pressed portion of the process cartridge is pressed by the first press portion of the image forming apparatus and positioned in place, whereby the process cartridge can be securely mounted in the predetermined position even in an incompletely mounted condition.\nIt is preferable that the engagement portion of the image forming apparatus comprises a recess or projection, and wherein the positioning member comprises a projection or recess which engages with the recess or projection of the engagement portion.\nPreferably, the first pressed portion is provided on the positioning member.\npreferably, the process cartridge further comprises a second pressed portion which is pressed by a second press portion provided on a cover of the image forming apparatus when closing the cover, the process cartridge is positioned at the first predetermined position by pressing the second press portion on the second pressed portion. In this case, the second pressed portion may be provided on the positioning member. The second pressed portion may be the first pressed portion.\nPreferably, the process cartridge may further comprises a third pressed portion which is pressed by a third press portion provided on the image forming apparatus when the process cartridge is moved from the second predetermined position to a third predetermined position to position the process cartridge at the third predetermined position. In this case, the pressing direction by the third press portion may be substantially perpendicular to the pressing direction by the first press portion. When the process cartridge is positioned at the third predetermined position, the first and second pressed portions are not pressed by the first and second press portion.\nPreferably, the pressing direction by the first and second press portions may be identical to the direction that the process cartridge is mounted on the image forming apparatus.\nAccording to a third aspect of the present invention, there is provided an image forming apparatus, comprising:\na process cartridge comprising:\na unit comprising an element used for an image forming process;\na positioning member having a positioning member;\na holding member for holding the unit to allow displacement of the unit relative to the positioning member; and\na restricting and releasing mechanism for restricting or releasing the displacement of the unit relative to the positioning member;\nan image forming unit for forming an image using the element of the process cartridge; and\nan engagement portion which engages with the positioning member of the process cartridge to position the process cartridge at a predetermined mount position in the image forming apparatus.\nAccording to a fourth aspect of the present invention, there is provided an image forming apparatus, comprising:\na process cartridge comprising:\na unit comprising an element used for an image forming process;\na positioning member having a first positioning member; and\na first pressed portion;\nan image forming unit for forming an image using the element of the process cartridge;\nan engagement portion which engages with the positioning member of the process cartridge to position the process cartridge at a first predetermined position in the image forming apparatus; and\na first press portion for pressing the first pressed portion of the process cartridge when the process cartridge is moved from the first predetermined position to a second predetermined position to position the process cartridge at the second predetermined position."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus used to hold work pieces for machining and, more specifically, to a vise capable of holding two work pieces.\n2. Description of the Prior Art\nNumerical control equipment are devices which perform functions automatically upon computer-given instructions. A single tool is able to perform a function and to move automatically among several work pieces that are secured to the work table of the device. The only function that is performed manually is securing the raw stock to the work table and removing the pieces upon completion of the machine's operations.\nThe operation of machining materials using numerically controlled equipment is much more expedient and efficient if more than one work piece, either raw stock or a partially finished piece, is available to the numerically controlled equipment at one time. Various dual vises are available which are capable of holding two work pieces. Whether the two pieces are at the same stage of production or at different stages, the ability to hold two pieces with one vise reduces the time needed to set up before an operation is performed by the numerically controlled equipment.\nA number of vises capable of holding two work pieces have been manufactured, with varying degrees of success. In U.S. Pat. No. 4,448,406, the moveable jaws are simultaneously driven in opposed reciprocation by a single shaft, the pressure exerted on both work pieces being the same.\nIn U.S. Pat. No. 4,341,375 a vise was designed for clamping a pair of snow skis. The vise uses a single shaft and is capable of clamping two skis of different widths. After the first ski is secured by the first vise assembly in a spring-tight position, continued rotation of the shaft causes pressure to be exerted on the second ski in the second vise assembly. After both skis are held spring-tight, further shaft rotation tightens each one equally. The operator is unable to apply varied pressures to the skis held in each vise assembly, and material of a delicate nature placed in the first vise assembly may be damaged by excessive force. A further problem with this design is that a ski held in the first vise assembly may slip from its original position while the second ski is being secured in the second vise assembly.\nOther improvements have been made to dual vises. The dual vise in U.S. Pat. No. 4,529,183 is capable of holding delicate material in both vise assemblies while holding the shaft in tension. The pressure exerted on each work piece in the vise assembly is equal. The operator has no option to differ the amount of pressure exerted on each work piece.\nWhile a work piece is being worked upon it needs to be securely held. However, the pressure should not be so excessive as to damage the finish of the work piece, or the work piece itself. A problem common to the prior art is an inability to hold two work pieces which may be at different stages of production without damaging the work piece at a more delicate stage, and yet to hold each work piece securely enough to avoid slippage."}
{"text": "Often, the most effective way to convey information to a human being is visually. Accordingly, millions of people work with a wide range of visual items in order to convey or receive information, and in order to collaborate. Such visual items might include, for example, concept sketches, engineering drawings, explosions of bills of materials, three-dimensional models depicting various structures such as buildings or molecular structures, training materials, illustrated installation instructions, planning diagrams, and so on.\nMore recently, these visual items are constructed electronically using, for example, Computer Aided Design (CAD) and solid modeling applications. Often these applications allow authors to attach data and constraints to the geometry. For instance, the application for constructing a bill of materials might allow for attributes such as part number and supplier to be associated with each part, the maximum angle between two components, or the like. An application that constructs an electronic version of an arena might have a tool for specifying a minimum clearance between seats, and so on.\nSuch applications have contributed enormously to the advancement of design and technology. However, any given application does have limits on the type of information that can be visually conveyed, how that information is visually conveyed, or the scope of data and behavior that can be attributed to the various visual representations. If the application is to be modified to go beyond these limits, a new application would typically be authored by a computer programmer which expands the capabilities of the application, or provides an entirely new application. Also, there are limits to how much a user (other than the actual author of the model) can manipulate the model to test various scenarios."}
{"text": "1. Field of the Invention\nThis invention relates to an authentication apparatus used for authenticating the user of an information processing machine, etc.\n2. Description of the Related Art\nIn recent years, the unauthorized use of an information processing machine has been socially at stake, and the management burden of personal information concerning authentication such as the user name and the password has grown. Against this backdrop, patent document 1 discloses an art of using an image processing technology to authenticate an individual based on the face image in a copier and permitting only the registered user to use the machine. According to the art, information that can be stolen such as the user name and the password is not used and thus the art has the advantage that it makes it possible to reduce the information management burden, etc. (Please see JP-A-2002-111924)"}
{"text": "All primary headaches are in need of better treatments. The most important primary headaches (i.e. independent disorders that are not caused by another disease) are migraine, tension-type headache and cluster headache. Migraine has a prevalence of 10% in the general population with a lifetime prevalence of 13% in men and 33% in women. Migraine is a highly disabling disease with high personal and social costs. To date, the precise mechanisms underlying the pathophysiology of migraine have remained elusive. Migraine strikes people during what are expected to be their most productive years: between ages 20 and 40 for most women, with a slightly higher age range for men. Migraine is typically characterized by unilateral onset of head pain, severe progressive intensity of pain, throbbing or pounding, and interference with the person's routine activities. Accompanying symptoms of photophobia (sensitivity to light) or phonosensitivity (intolerance to noise), as well as nausea and/or vomiting, are common, and often leads to the inability to perform daily tasks. A large portion of people with migraine often have no accompanying pain, their predominant symptom instead being vertigo (a spinning sensation) or dizziness/disequilibrium (balance loss), mental confusion, disorientation, dysarthria, visual distortion or altered visual clarity, or extremity paresis. Patients with migraine associated vertigo (MAV) are often seen by audiologists and vestibular rehabilitation therapists for evaluation and treatment. Because the exact mechanisms of migraine are still not completely understood, the management of migraine dizziness presently includes a combination of medications, vestibular rehabilitation, and lifestyle modifications that include limitation of risk factors associated with migraine (those related to diet, sleep, stress, exercise, and environmental factors).\nMigraine is a disease associated with increased synthesis and release of calcitonin gene related peptide (CGRP) and a migraine attack can be blocked with CGRP antagonists. The actual pain is generated by nociceptors of trigeminal nerve endings in the dura. Low serotonin levels may sensitize the nociceptors of trigeminal neurons. Triptans and ergotamins, which decrease serotonin, are associated with relief of acute pain. In contrast, tricyclic antidepressants and selective serotonin and noradrenaline reuptake inhibitors, which are associated with increases in serotonin, are utilized for migraine prevention.\nMigraine attacks can be triggered by intrinsic cerebral factors (e.g. calcitonin gene related peptide (CGRP) release), nitric oxide like tri-nitroglycerine, corticotrophin releasing hormone (stress), pro-inflammatory cytokines, and degranulation of mast cells located in the dura. While migraine has a genetic background, twin studies reveal that the cause of a majority of migraines appears to be due to environmental factors. The mechanism of triggering migraine is, however, still not understood. The cause of both migraines and chronic dizziness has eluded investigators for centuries and it therefore presents a truly long felt but unsolved mystery as to its causation and treatment."}
{"text": "1. Field of the Invention\nThe present invention relates to a drive unit, and, more particularly, to a drive unit for a motor vehicle having an internal combustion engine with a cooling system having a fan and a hydrodynamic retarder having a stator, a rotor and a housing. Such a drive unit is known from DE 37 13 580 C1.\n2. Description of the Related Art\nRetarders are primarily employed in heavy vehicles to absorb and convert to heat the kinetic braking energy which accrues most notably during high speed braking actions (adaptation braking). Retarders are also well suited for situations involving sustained braking, for example, maintaining a constant speed of 30 km/h on an incline of 7%. Oil normally serves as the operating fluid. The heat transferred in the retarder to the operating fluid must be delivered, by means of a specific heat exchanger, to a coolant or the ambient air.\nU.S. Pat. No. 3,720,372 discloses a retarder which in one embodiment is powered by the crankshaft and constantly flooded by the coolant of the cooling system. The rotor of the retarder serves as circulating pump directing heated coolant to a heater core within the passenger compartment of the automobile. The purpose of this system is to heat the coolant by means of the retarder and thereby heat the passenger compartment. There is also a control system arranged on the retarder which controls the distribution of the coolant depending on its temperature.\nAlso known, from DE-PS 33 01 560, is a retarder which, by way of a clutch, is connected to the crankshaft of the drive engine and to the driven wheels of the vehicle. It is not the purpose of this retarder, however, to absorb the high kinetic braking energy of the vehicle and convert it to heat. Instead, the retarder is operated exclusively as a heater, with the available operating energy input controlling the heating output. The coolant of the engine also serves as the operating fluid of the retarder.\nA retarder known from DE-AS 1 946 167 (U.S. Pat. No. 3,650,358) is powered by the crankshaft of an internal combustion engine whose coolant also serves as the operating fluid for the retarder. The advantage of this mode of operation is that the accruing heat develops directly in the coolant passed to the radiator and a heat exchanger between the two fluids is unnecessary. The rotor is mounted on an antifriction bearing and the seal between the frame and rotor shaft is established by two lip seals.\nIt is desirable to keep the overall axial dimensions and weight of the drive units of this type as low as possible, especially when the drive unit will be used in a motor vehicle. The drive units known heretofore have failed to satisfy this goal."}
{"text": "The present invention relates to measuring and reporting link quality in a cellular telecommunication system. More particularly, the present invention relates to using frame erasure rates (also sometimes referred to as xe2x80x9cframe error ratesxe2x80x9d) in link quality measurement to improve network quality observation.\nIn most radiocommunication systems, such as cellular telephone systems, one or both of the base station and the mobile unit typically evaluate and report on the signal quality of the communication link (herein referred to as link quality) between the mobile and the base station. In older analog systems such as the Advanced Mobile Phone Service (AMPS) and Extended Total Access Communication System (ETACS), mobile units participate in the evaluation of link quality by transmitting a pilot tone to the base station. The base station then measures the signal strength of the pilot tone to estimate uplink and downlink signal quality.\nIn more contemporary digital systems such as the Global System for Mobile communication (GSM), used as a standard in Europe, and the Digital Advanced Mobile Phone Service (D-AMPS), used as a standard in North America, the mobile units measure the downlink quality and then report the measurements back to the base station. Typically, mobile units measure and report the bit error rate (BER). In addition to measuring and reporting the signal quality of the downlink between the base station and the mobile unit, each mobile unit is also capable of measuring and reporting the received signal strength (RSS) for a number of alternative frequency channels herein referred to as candidate channels. In a similar manner, the base station may measure the received signal strength and BER for the uplink. Both uplink and downlink quality information can then be reported to the network for storage and analysis by the operator. If necessary, remedial action can be taken at the base station to improve a connection\"\"s signal quality.\nFor example, if the BER for the frequency channel currently being used for the downlink is greater than desirable and the corresponding signal strength is less than desirable, the base station may take remedial action. The base station may exchange the current frequency channel with one of the candidate channels or may even decide to xe2x80x9chandoffxe2x80x9d the mobile unit to a neighboring base station. These remedial actions help to insure that the user of the mobile unit experiences the best possible signal quality.\nTo avoid taking unnecessary remedial action in response to a few poor link quality measurements, existing systems average the link quality measurements in accordance with a time constant. The time constant will, in general, vary from a few milliseconds to several seconds depending upon the data rate, the communication protocol, and the frequency of the communication channel. For example, systems operating in accordance with GSM estimate the BER for each 20 ms speech frame and average the estimated BERs over half second intervals. Unfortunately, averaging the link quality measurements made using BER information tends to present an inaccurate measure of link quality. This is especially true when the operating environment is subjected to the effects of Rayleigh fading or other time dependent interference which increase the peak-to-average ratio of the signal quality measurements by approximately 3 to 4 decibels, i.e., degrade link quality without significantly affecting the average BER measurement.\nThis, in turn, makes it difficult for network operators to accurately determine the actual quality perceived by users in different parts of the network. Applicants have recognized that the problems underlying these conventional link quality reporting techniques involve the usage of RSS and/or BER information to estimate perceived link quality. In particular, the correlation between these parameters and perceived signal quality may be low and, moreover, the quality of the BER estimate itself may be poor.\nAccordingly, it would be desirable to provide new systems and techniques for estimating and reporting link quality in a radiocommunication system.\nExemplary embodiments of the present invention provide techniques and systems for measuring and reporting frame erasures in either or both the uplink and downlink of communication channels. This information is used by the network operator to more accurately estimate perceived signal quality in the system.\nAccording to some exemplary embodiments, the frame erasure rate is measured and reported. According to other exemplary embodiments, the time distribution of the frame erasures is also reported.\nVarious processing techniques are used according to the present invention to present the network operator with the link quality information needed to optimize system parameters. For example, frame erasure information can be presented on a per connection, per cell or system-wide basis. Moreover, the link quality information associated with connections that have recently been handed-off can be highlighted to alert network operators to potential problems in a particular area\"\"s handoff algorithm. If the data reported by the system to the network operator is at a higher level than the desired presentation of that data, then exemplary embodiments of the present invention provide for transforming the received data to provide an estimate of the lower level presentation data."}
{"text": "The present invention relates to a propylene-ethylene block copolymer composition.\nPropylene-ethylene block copolymers are used as sheets, molded articles, and also as materials for production of car parts, home electric appliances, and so forth because they are superior in stiffness and low temperature impact strength. Because of poor transparency, however, they can find only limited applications. Thus, to increase the transparency of the copolymers, it has been proposed to add nucleating agents. In this case, a problem arises in that the low temperature impact strength drops. Furthermore, an attempt has been made to increase the low temperature impact strength of propylene-ethylene block copolymers in order to extend their applications by compounding inorganic filler such as talc. Addition of such inorganic filler, however, results in a decrease in transparency of the copolymers.\nThus it has been desired to improve the transparency of propylene-ethylene block copolymers without deteriorating their superior stiffness and low temperature impact strength or while at the same time more increasing the stiffness and low temperature impact strength."}
{"text": "Eribulin (marketed under the trade name HALAVEN® as eribulin mesylate), a nontaxane microtubule dynamics inhibitor, is a structurally simplified, synthetic analog of the marine natural product halichondrin B. Methods for the synthesis of eribulin and other halichondrin B analogs are described in U.S. Pat. Nos. 6,214,865, 6,365,759, 6,469,182, 7,982,060, and 8,148,554, the syntheses of which are incorporated herein by reference. New methods for the synthesis of halichondrin B analogs, in particular eribulin and eribulin mesylate, are desirable."}
{"text": "Clinical studies and practice have shown that providing a reduced pressure in proximity to a tissue site augments and accelerates the growth of new tissue at the tissue site. The applications of this phenomenon are numerous, but application of reduced pressure has been particularly successful in treating wounds. This treatment (frequently referred to in the medical community as “negative pressure wound therapy,” “reduced pressure therapy,” or “vacuum therapy”) provides a number of benefits, including faster healing, and increased formulation of granulation tissue."}
{"text": "1. Field of the Invention\nThis invention relates to a device for collecting solar heat by the use of a vacuum tube type heat collector having heat collecting plates and heat medium tubes.\n2. Prior Art\nThis type of device is designed to accumulate a high-temperature heat medium in a heat accumulating tank by repeated operation of transferring a heat medium heated to high temperatures by a vacuum tube type heat collector to a heat accumulating tank and of feeding a low-temperature heat medium in the heat accumulating tank to a heat collector. Generally, the heat medium is circulated by a circulation pump between the heat collector and the heat accumulating tank, and the pump is driven when the temperature of the heat medium heated at the outlet portion of the heat collector is raised higher than a specified temperature or when the temperature on the side of the heat collector is higher by comparison between the heat medium temperature at the bottom of the tank and that at the outlet portion of the heat collector.\nThe vacuum tube type heat collector provides the advantage that the collector permits the free installation thereof in various ways such as horizontal installation, but it has the counterbalancing disadvantage that the point of space wherein the temperature is highest inside the heat collector is difficult of easy measurement as opposed to a plate type heat collector.\nUnder examples of heat controlling devices for a conventional vacuum tube type heat collector are included a device for controlling heat collection by an insolation sensor, a device for controlling heat collection by the temperature of a header portion (confluence of pipes through which a heat medium passes), and the like.\nFIG. 1 shows a conventional type device for controlling the collection of heat by an insolation sensor. A low-temperature sensor 2 is disposed in the heat accumulating tank 1, and the insolation sensor 4 is installed in matching relation with respect to a degree of inclination of a heat collecting plate 3. The insolation is sensed by the sensor 4 and is compared by a heat collection controlling unit 5 with the low-temperature sensor 2 on the bottom side of the heat accumulating tank 1, and when it is approximately judged by the controlling unit 5 that collection of heat is possible, a pump 6 operates to collect heat.\nA device for controlling the collection of heat by the temperature of the header portion is of the construction in which the header is provided with a sensor so as to sense the temperature of the header and that of the heat medium inside the header and in which when the temperatures of the header and the heat medium in the header become higher than the temperature of the heat accumulating tank to a certain set value, the heat collection pump operates and when the temperatures become lower than the set value, the pump ceases to operate.\nFIG. 2 shows a flow chart representing the control of the above-mentioned device using an insolation sensor. In the figure, the reference character TL designates a temperature in the heat accumulating tank 1 detected by the low-temperature sensor 2, and I designates an insolation detected by an insolation sensor 4. FIG. 3 shows a relation between the insolation I and the temperature TL in the heat accumulating tank 1. Both of the aforementioned devices relate to a method of sensing the rough temperature in the insolation sensor and the header portion respectively, and, when the insolation sensor is used, the relation between the insolation and the temperature describes a secondary curve as shown in FIG. 3, and the characters I and TL are not brought into a relation of linear proportion. When the temperature of the header portion is sensed, transfer of heat from the heat collecting unit is delayed, thus the devices being reduced in performance, and positive control becomes impossible. When the insolation is reduced, or when the weather conditions change, or when the temperature of the open air is reduced, the devices are subject to external conditions to render it impossible to make accurate control."}
{"text": "Aircraft wheels may use tapered roller bearings to enable wheels to roll smoothly. The typical wheel bearing may have a life span that is shorter than the life of the aircraft. Thus, bearings may be replaced from time to time. Visual inspections may be performed to evaluate whether a bearing may be ready for replacement. However, complete replacement of the bearing may cause slight damage to the wheel as a portion of the bearing may be press fit into the wheel. Thus, after a limited number of bearing replacement, the damage to the wheel may dictate that the wheel be replaced as well.\nOne obstacle to preventative bearing maintenance is human error and consistency. Visual inspection quality may vary by the inspector or process used. Additionally, some replacement bearings may be misaligned upon reassembly. If misalignment goes undetected, a bearing failure may result. Bearings may fail due to corrosion, spalling, misalignment, foreign materials, handling damage, brinelling, lubrication, or excessive load. If bearing failure does occur, it also may result in permanent damage to the wheel or axle. In some instances, wheels may detach from the axle and aircraft. In the field, undetected wheel-bearing failures frequently dictate a high-cost replacement of the wheel and axle."}
{"text": "Field of Invention\nThis disclosure relates to an image sensor, and specifically to an image sensor having an extended spectral range.\nDescription of Related Technology\nImage processing devices, such as digital cameras, smartphones or tablet computers, rely on image sensors to capture images. Image sensors receive light and convert that light into electrical signals. The image processing devices then transform these electrical signals into digital images.\nDifferent types of image sensors are currently available. For example, image processing devices typically utilize either a frontside-illumination (FSI) image sensor or a backside-illumination (BSI) image sensor. An FSI image sensor is typically oriented such that light enters the top of the FSI image sensor and passes through a metal-interconnect layer before striking the light-sensing surface. In contrast, BSI image sensors are oriented to allow light to enter from the top of the BSI image sensor and to strike a light-sensing surface without passing through the metal/wiring layer of the BSI wafer. While each of the FSI and BSI image sensors have favorable imaging characteristics, they both may have limited spectral responses. Accordingly, there is a need for an image sensor having a greater spectral response than either a FSI or BSI image sensor."}
{"text": "1. Technical Field\nThe disclosure relates to a charging battery unit and an electronic device with the same.\n2. Description of Related Art\nWireless charging device used in the portable electronic product needs to pass through Wireless Power Consortium (WPC, Qi) authentication. However, the receiving unit of the wireless charging device is incorporated into the back cover of the portable electronic product; therefore, it needs to detect the whole portable electronic product in the WPC authentication, not only the receiving unit, resulting in the detecting process inconvenient.\nTherefore, what is needed is a charging battery unit and an electronic device with the same which can overcome the described limitations."}
{"text": "This description relates to operation of security and surveillance systems.\nIt is common for businesses and homeowners to have a security system for detecting alarm conditions at their premises and signaling the conditions to a monitoring station or to authorized users of the security system. Security systems often include an intrusion detection panel that is electrically or wirelessly connected to a variety of sensors. Those sensors types typically include motion detectors, cameras, and proximity sensors (used to determine whether a door or window has been opened)."}
{"text": "The present invention relates to the synthesis of conjugated polyene sterol derivatives, the compounds obtained and to their use as fluorescent probes for cellular membranes.\nA membrane probe is essentially a research tool used to probe (investigate) the environment of a membrane. Membrane probes find application, for example, in determining the effect of environmental pollutants and dietary constituents on cell membranes. Cell membranes are very selective and will permit only selected molecules to permeate the membrane. As well, the interactions between the membrane and its environment (i.e. the molecules in the environment) are very selective. Membrane probes are useful in the investigation of the structure and function of membrane constituents such as proteins and enzymes. Cholesterol is a major membrane constituent. The molecular structure of the probes of the present invention is very similar to cholesterol and thus the probes behave in a manner similar to cholesterol.\nThe probes of the present invention are fluorescent, so their behaviour in the membrane environment and their effect on the membrane can be easily determined by simply shining light on the membrane and observing the light emitted by the probe. Because of the very great sensitivity of fluorescence techniques, only very small amounts of extrinsic probe material need be incorporated into the sample of interest. This contrasts with other spectroscopic techniques such as nuclear magnetic resonance, absorption and electron spin resonance spectroscopy, which require significantly greater amounts of probe substance, with much greater risk of altering the system being investigated. Techniques using radio-isotopes have a high degree of sensitivity but by their nature are more hazardous and the radioisotopic probe materials have limited shelf life.\nBy studying the behiavour of the probe of the invention more can be learned about the behaviour of cholesterol. The probes of the present invention have potential for use in determining cholesterol levels and cholesterol properties in membranes and cell membrane properties and can be applied to clinical assays and diagnoses involving cholesterol.\nAromatic olefins, particularly diphenylhexatriene (DPH), have been used extensively as fluorescent probes of membrane fluidity, (see for example, L. A. Chen, R. E. Dale, S. Roth and L. Brand, J. Biol. Chem., 1977, 252, 2163). Such molecules, however, are not natural membrane constituents and can provide only indirect insight into protein-lipid and lipid-lipid interactions.\nCholesterol is an important lipid component, but its role in influencing membrane structure and dynamics is poorly understood. Some cholesterol-type molecules in which the C-3 alcohol function has been modified have been synthesized. (See for example, R. R. Rando, F. W. Bangerter and M. R. Alecio, Biochim. Biophys. Acta, 1981, 684, 12.) For example, the compounds ##STR1##\nApart from a loss of the amphipathic nature of the cholesterol moiety in some of these derivatives (e.g. 1(a)), in all cases the C-3 hydroxy function has been drastically altered. Other membrane probes have been synthesized in which the C-3 hydroxy group is intact. (See for example, R. Bergeron and J. Scott, Anal. Biochem., 1982, 119, 12H; F. Schroeder, FEBS Lett., 1981, 1335, 127.) For example, the compounds ##STR2## However, the unsaturation in the ring systems 2(a) and 2(b) drastically changes the molecular geometry. Such probes would be expected to pack differently from cholesterol in membranes and the lipid-probe interactions would be different from those of lipid-cholesterol.\nFluorescent chromophores in the C-17 side-chain of cholesterol have been reported. (See for example, Y. J. Kao, A. K. Soutar, K. -Y. Hong, H. J. Pownall and L. C. Smith, Biochemistry, 1978, 17, 2689.) For example, the compounds ##STR3##\nIn compounds of structures 3(a), 3(b) and 3(c) the hydrophobic character of the C-17 side-chain is much different from that of cholesterol. The membrane probes of the present invention, which have an olefinic fluorescent chromophore in the C-17 side-chain, resemble cholesterol more closely both in geometry and in amphipathic nature."}
{"text": "The invention relates to a method suitable for the oxidation of a conjugated diolefin. In another aspect the invention relates to a composition useful as a catalyst.\nIt is desirable to oxidize conjugated diolefins, such as 1,3-butadiene and/or 2-methyl-1,3-butadiene to produce various compounds such as the ethylenically unsaturated esters of such diolefins. A more specific illustration is the oxidation of 1,3-butadiene to produce 1,4-diacetoxy-2-butene. The diacetoxybutene is then easily converted, by processes well known in the art, to other compounds such as tetrahydrofuran or 1,4-butanediol. Although there are various processes and catalysts known which are useful for the oxidation of a conjugated diolefin, most of these processes are relatively expensive to carry out and frequently corrosion of process equipment is a problem. Therefore new processes and catalysts are desirable in an effort to more fully develop the art and improve the overall process.\nAn object of the present invention is to oxidize a conjugated diolefin.\nAnother object of the invention is to oxidize a conjugated diolefin more economically than can be done presently.\nAnother object of the invention is to provide a catalyst useful for the oxidation of conjugated diolefins.\nOther objects, advantages and aspects of the present invention will be apparent to those skilled in the art after studying the specification and appended claims."}
{"text": "Numerous diseases are associated with inherited or somatic mutations. In many cases, these mutations are present in the transcript region of genes, the products of which control important physiological functions including, for example, gene expression, cell signaling, tissue structure, and the metabolism and catabolism of various biomolecules. Mutations in these genes, which are often only single nucleotide changes (e.g., non-sense mutations, missense mutations), can have negative effects on the expression, stability and/or function of the gene product resulting in alterations in one or more physiological functions.\nA number of different mutations have been identified in the Alpha-1 antitrypsin (AAT) gene. AAT is one of the primary circulating serum anti-proteases in humans. AAT inhibits a variety of serine proteinases, with neutrophil elastase being one of the most physiologically important, as well as inhibiting a number of metalloproteinases and other pro-inflammatory and pro-apoptotic molecules. AAT is normally produced within hepatocytes and macrophages, where hepatocyte-derived AAT forms the bulk of the physiologic reserve of AAT.\nApproximately 4% of the North American and Northern European populations possess at least one copy of a mutant allele, known as PI*Z (Z-AAT) which results from a single amino acid substitution of lysine for glutamate at position 342 in the mature protein (position 366 in the precursor protein). In the homozygous state, this mutation leads to severe deficiency of AAT, and can result in two distinct pathologic states: a lung disease which is primarily due to the loss of antiprotease function, and a liver disease (present to a significant degree in approximately 10-15% of patients) due to a toxic gain of function of the Z-AAT mutant protein.\nInvestigational clinical gene therapy products for gene augmentation of AAT have been developed as potential treatments for lung disease using the recombinant adeno-associated viral (rAAV) vectors. Researchers have also applied genetic technologies in an effort to down-regulate the levels of AAT mRNA. One approach was to utilize hammerhead ribozymes designed to cleave AAT mRNA at a specific site. Another approach involves the use of RNA interference to decrease levels of the mutant mRNA transcript."}
{"text": "1. Field of the Invention\nThe present invention generally relates to anti-decoupling arrangements for connectors of the type in which coupling is achieved by means of a coupling nut, and more particularly to an anti-decoupling arrangement for an electrical connector that uses a ratchet mechanism to limit rotation of the coupling nut in the decoupling direction and a spiral lock clutch to permit free rotation of the coupling nut in the coupling direction. Still more particularly, the invention relates to improvements on the anti-decoupling arrangement disclosed in copending U.S. patent application Ser. No. 09/391,458, filed Sep. 8, 1999, incorporated by reference herein.\n2. Description of the Related Art\nA typical connector to which the present invention may be applied includes a connector shell containing electrical contacts and an internally threaded coupling nut rotatably mounted on the connector shell. The connector shell is coupled to a corresponding externally threaded mating connector by means of the coupling nut in such a manner that electrical contacts in the mating connector engage the electrical contacts in the connector shell. The coupling nut is held on the connector shell by one or more retaining rings and/or spring washers that are designed to captivate or press a radial flange of the coupling nut against a corresponding flange or shoulder on the connector shell.\nBecause the frictional anti-locking force generated by engagement between the coupling nut and connector shell in such an arrangement is insufficient to prevent the coupling nut from rotating in a decoupling direction as a result of vibrations or shocks, compromising seals and possibly affecting the integrity of the electrical connections between contacts, it is conventional to include an additional anti-decoupling mechanism in connectors likely to be used in environments where vibrations or shocks are likely to occur, such as in military high-performance aircraft and other vehicles. The simplest and most common method of preventing unintended decoupling as a result of shocks or vibrations has been to include in the connector a metal ratchet spring having protrusions or dimples at the center of the beam, the ratchet spring being permanently attached to the inside diameter of the threaded coupling nut. The connector shell is provided with ratchet teeth on its outer diameter, which are engaged by the ratchet spring.\nOne problem with this type of coupling is that the discrete detent positions do not necessarily lie in phase with the fully clamped position of the ring, such that even slight vibrations can cause the ring to back off slightly, which can cause sealing problems. In addition, the detent members in this configuration have very little effective surface area, causing rapid wearing away of the teeth on the ratchet wheel each time the connector is mated or unmated.\nA solution to the problems of wear and phasing of the ratchet teeth and detents is described in copending U.S. patent application Ser. No. 09/391,458, which is directed to various improvements in a spiral lock clutch anti-decoupling mechanism originally proposed in U.S. Pat. No. 4,536,048. The anti-decoupling mechanism described in the copending patent application includes a spiral lock clutch that permits free running in the coupling direction, a spring ring, and a tooth wheel all surrounding a connector shell and captured between a snap-ring on the connector shell and an inwardly extending flange on the coupling nut. The tooth wheel includes extensions or knurls that cooperate with corresponding slots or surfaces of the coupling nut to prevent relative rotation between the coupling nut and the tooth wheel, while the spring ring includes spring tines that engage radial cuts in the tooth wheel to permit ratcheting of the tooth wheel relative to the spring ring. The spring ring, in turn, is locked against rotation relative to the spiral lock clutch. During coupling, turning of the coupling nut causes corresponding turning of the tooth wheel. Since the spiral lock clutch is arranged to unwind and permit free running in the coupling direction, the engagement between the spring tines on the spring ring and the radial cuts is not subject to any ratcheting force and the spring ring and spiral lock clutch turn freely with the coupling nut and tooth wheel. During uncoupling, on the other hand, the spiral lock clutch winds tightly against the connector shell, preventing rotation of the spiral lock clutch and spring ring. In order to permit the coupling nut to rotate, a sufficient force must be applied to the coupling nut to permit ratcheting of the spring ring relative to the tooth wheel, i.e., to permit the spring tines to glide over the teeth formed by the radial cuts in the ratchet wheel.\nThe above-described anti-decoupling arrangement has the advantages, relative to the anti-decoupling arrangement described in U.S. Pat. No. 4,536,048, of attaining a high uncoupling torque due to the use of multiple tines or beams on the spring ring attached to the spiral lock clutch, control of the coupling torque through appropriate choice of the spiral lock clutch, spring tines, and tooth configuration, and simplified assembly to the connector shell by fitting all of the components over the shell, angularly orienting the components, and holding them in place with a retaining ring. Nevertheless, the above-described anti-decoupling mechanism still could benefit from the following improvements:\n(i) a greater degree of adjustment of the de-coupling torque; PA1 (ii) a still higher de-coupling torque than can be achieved with the prior arrangement; PA1 (iii) smoother non-binding operation; and PA1 (iv) a less critical assembly method. PA1 (i) arranging the ratchet assembly cantilever beams so that they operate radially outwardly rather than axially; and PA1 (ii) eliminating the ratchet assembly detent ring (i.e., the toothed wheel) used in the prior anti-decoupling arrangement in favor of serrations formed into the inner diameter of the coupling nut.\nThese improvements are achieved by modifying the anti-decoupling device described in the copending patent application so that the clutch mechanism and the ratchet mechanism operate completely independently of one another in a non-interfering manner, and in particular by:\nThese modifications not only reduce the number of components and also provide mechanical advantages that increase the range of possible decoupling torques, but they also eliminate any interference between the coupling nut shoulder and the back side of the spiral wound clutch band so as to provide a smoother coupling feel and a more positive and stronger clutch grip, eliminate press fits or keyed components that complicate assembly, permit a stronger and more easily assembled attachment of the spring ring to the spiral wound clutch, reduce tolerance build-up between components (due to the smaller number of axially stacked components), and make it possible to more easily disassemble the anti-coupling mechanism for repair or torque adjustment."}
{"text": "This application is based on Patent Application Nos. 92120/1998 filed on Apr. 3, 1998 in Japan and 205706/1998 filed on Jul. 21, 1998 in Japan, the content of which is incorporated hereinto by reference.\n1. Field of the Invention\nThe invention relates to a method for adjusting dot forming or depositing positions in dot matrix recording and a printing apparatus using the method. More particularly, the invention relates to a method for adjusting dot forming positions, which are applicable to printing registration in the case of bi-directionally printing by a forward and reverse scan of a print head or to printing registration in the case of printing by means of a plurality of print heads, and printing apparatus using the method.\n2. Description of the Related Art\nIn recent years, the office automation instruments such as the personal computer and the word processor which is relatively cheap are widely used, and an improvement in high-speed technique and an improvement in high image quality technique of various recording apparatuses for printing-out the information which are entered by the instruments are developed rapidly. In recording apparatuses, a serial printer using a dot matrix recording (printing) method comes to attention as a recording apparatus (a printing apparatus) which realizes printing of a high speed or high image quality with the low cost. For such printers, as the technique which prints at high speed, for example there is a bi-directional printing method and as the technique which the prints in high image quality, for example, there is a multi scanning printing method.\n[Bi-directional Printing Method]\nAs the improvement in high-speed technique, in a printing head which has a plurality of printing elements, although it is also thought to plan an increase in the number of a printing elements and an improvement in a scanning speed of the print head, it is also an effective method to perform bi-directional printing scannings of the print head.\nAlthough, since there is usually the time required for paper-feeding and paper-discharging or the like, it does not become a simply proportional relation, in the bi-directional printing a printing speed of approximately two times can be obtained as compared with the one-directional printing in the printing apparatus.\nFor example, when using the print head which the 64 pieces of ejection openings are arranged with 360 dpi (dots/inch) in printing density in a direction different from the printing scanning (main scanning) direction (for example, in a sub-scanning direction which is the feeding direction of the printing medium), a printing is performed on, a printing medium of A4 size set in the direction of the length, the printing can be completed by scanning of approximately 60 times. The reason is that, in one-directional printing, each printing scanning is performed only at the time of the movement in the one direction from the predetermined scanning commencement position, and since non-printing scanning to the inverse direction for returning to the scanning commencement position from a scanning completion position is attended, reciprocation of approximately 60 times is required. On the other hand, printing is completed by the reciprocating printing scanning of approximately 30 times in bi-directional printing, so that printing can be performed and since it becomes possible on at the speed of approximately 2 times, whereby bi-directional printing can be considered to be an effective method for an improvement in a printing speed.\nIn order to register dot-forming positions (for example, for an ink jet printing apparatus, a deposition or landing position of ink) at a forward trip and a return trip together in such bi-directional printing, using a position detection means such as an encoder, based on the detecting position, printing timing is controlled. However, it has been thought that since to form such a feedback controlled system causes an increase in the cost of the printing apparatus, it is difficult to realize this, in the printing apparatus which is relatively cheap.\n(Multi Scanning Printing Method)\nSecondly, a multi scanning printing method is explained as one example of the improvement in high image quality technique.\nWhen printing is performed using the print head which has a plurality of printing elements, quality of the printed image depends on performance of a print head itself greatly. For example, in the case of the ink-jet print head, the slight differences, which is generated in a print head manufacturing step, such as variations of a form of ink ejection openings and the elements for generating energy for ejecting ink such as an electro-thermal converting elements (ejection heaters), influence a direction and an amount of ejected ink, and result in the cause which makes the unevenness in density of the image which is formed finally to reduce the image quality.\nSpecific examples are described using FIGS. 1A to 1C and FIGS. 2A to 2C. Referring to FIG. 1A, a reference numeral 201 denotes a print head, and for simplicity, is constituted by the eight pieces of nozzles 202 (herein, as far as not mentioned specifically, refer to the ejection opening, the liquid passage communicated with this opening and the element for generating an energy used for ink, in summary). A reference numeral 203 denotes the ink, for example, which are ejected as a drop from the nozzle 202. It is ideal that the ink is ejected from each ejection opening by the approximately uniform amount of discharge and in the justified direction as shown in this drawings. When such discharge is performed, as shown in FIG. 1B, ink dots which are justified in size are deposited or landed on the printing medium and, as shown in FIG. 1C, the uniform images that there is no unevenness in density also as a whole can be obtained.\nHowever, there are the variations in the nozzles in the print head 20 actually as is mentioned above, and when printing is performed as mentioned above as it is, as shown in FIG. 2A, the variations are caused in size of the ink drops and in the ejecting direction of ink discharged from nozzles and the ink drops are deposit or landed on a printing medium as shown in FIG. 2B. In this drawing, a part of the white paper that an area factor can not be served up to 100% periodically exists with respect to the horizontal scanning direction of the head, moreover, in contrast with this, the dots are overlapped each other more than required or white stripes as shown in the center of this drawing have been generated. A gathering of the landed dots in such condition forms the density distribution shown in FIG. 2C to the direction in which nozzles are arranged, and the result is that, so far as usually seen by eyes of a human, these objects are sensed as the unevenness in density.\nTherefore, as a countermeasure of this unevenness in density, the following method has been devised. The method is described. using FIGS. 3A to 3C and FIGS. 4A to 4B.\nAccording to this method, in order that the printing with regard to the same region as shown in FIG. 1 and FIG. 2 is made to be completed, the print head 201 is scanned 3 times as shown in FIG. 3A and FIGS. 4A to 4C. The region defining four pixels which is a half of eight pixels as a unit in the direction of length in the drawing has been completed by two passes. In this case, the 8 nozzles of the print head are divided into a group of 4 nozzles of upper half and 4 nozzles of lower half in the drawing and the dots which one nozzle forms by scanning of one time are the dots that the image data are thinned into approximately a half in accordance with the certain predetermined image data arrangement. Moreover, at the second scanning, the dots are embedded in the image data of the half of the remaining and the regions defined four pixels as the unit are completed progressively. Hereinafter, the printing method described above is referred to as a multi scanning printing method.\nUsing such printing method, even when the print head 201 which is equal to the print head 201 shown in FIG. 2 are used, the influence to the printed image by the variations of each nozzle is reduced by half, whereby the printed image becomes as shown in FIG. 3B and no black stripe and white stripe as shown in FIG. 2B becomes easy to be seen. Therefore, the unevenness in density is fairly also mitigated as compared with the case of FIG. 2C as shown in FIG. 3C.\nWhen such printing is performed, although at first scanning and at second scanning, the image data are mutually divided in a manner to be complemental each other in accordance with the certain predetermined arrangement (a mask), usually, this image data arrangement (the thinned patterns) as shown in FIG. 4A to FIG. 4C, at every one pixel arranged in rows and columns, it is most general to use the formation which makes to form a checker or lattice matrix.\nIn a unit printing region (here, per four pixels), printing is completed by the first scanning which forms the dots into the checker or lattice pattern and the second scanning which forms the dots into the inverted checker or lattice pattern.\nMoreover, usually, travel (vertical scanning travel) of the printing medium between each main scanning is established at a constant, and in the case of FIG. 3 and FIG. 4, is made to move every four nozzles equally.\n(Dot Alignment)\nAs an example of the other improvement in high image quality technique in the dot matrix printing method, there is a dot alignment technique adjusting the dot depositing position. A dot alignment is an adjustment method adjusting the positions which the dots on the printing medium have formed by any means, and in general, the prior dot alignment has been performed as follows.\nFor example, a ruled line or the like is printed on a printing medium in depositing registration of the forward scan and the reverse scan upon reciprocal or bi-directional printing by adjusting printing timing in the forward scan and the reverse scan respectively, while a relative printing position condition in reciprocal scan is varied. The results of printing has been observed by a user oneself to select the printing condition where best printing registration is achieved, that is, the condition that printing is performed without offset of the ruled line or the like and to set the condition directly into the printing apparatus by entering through a key-operation or the like or to set the depositing position condition into the printing apparatus by operating a host computer through an application.\nMoreover, the ruled line or the like is printed on the medium under printing in the printing apparatus having a plurality of heads, when printing is performed between a plurality of heads, while a relative printing position condition between a plurality of heads is varied, with the respective head. As is mentioned above, the optimum condition that best printing registration is achieved has been selected to vary the relative printing position condition to set the printing position condition into the printing apparatus every each head in the mentioned-above manner.\nHere, the case where the offset of the dots has been occurred is described.\n(Problems Upon Performing Image-formation by Bi-directional Printing)\nDue to bi-directional printing, the following problems has been caused.\nFirst, when the ruled line (the ruled line of the longitudinal direction) in the direction perpendicular to the horizontal scan of the print head is printed, between the ruled line element which is printed in the forward scan and the ruled line element which is printed in the reverse scan, the dot depositing positions are not registered and the ruled line is not formed into a straight line, but a difference in level occurs. This is referred to as a so-called xe2x80x9coffset in ruled linexe2x80x9d, and this is considered to be the most general disorder which can be recognized by the usual users. In the many cases, the ruled line is formed by a black color, whereby, though the offset in ruled line has been understood as the problem where a monochrome image is formed generally, a similar phenomenon can be caused in the color image also.\nWhen the multi scanning printing is used along with bi-directional printing in order to improve in high image quality, even though in bi-directional printing the depositing positions are not registered, as an effect of the multi scanning printing the offset in the pixel level is not easy to be seen, but from a macroscopic viewpoint the entire image can be seen unequally and is recognized as an unpleasant figure by the user. This generally is called as a texture, and appears on the image in the specific period where there is the offset in the delicate depositing position, thereby being caused. In a strong image in contrast such as the monochrome it is easy to be seen, moreover, when for the printing medium capable of high-density printing such as a coat paper middle-tones printing is performed, it can be easy to be seen.\n(Problems in the Case of Performing the Image Formation Using a Plurality of the Print Heads)\nIn the printing apparatus having a plurality of heads, the problems of the case where the offset in the depositing positions of the dots between a plurality of heads has been occurred is discussed.\nWhen the image printing is performed, several colors are combined to perform the image formation frequently, and it is general to use four colors which added black in addition to three primary colors of yellow, magenta and cyan and it is used most abundantly. When in the case where a plurality of print heads for printing these colors are used, there is the offset of the depositing positions between the print heads, depending upon the amount of the offset, when a different color one another is about to be printed on the same pixel, a deviation in color matching is caused. For example, magenta and cyan are used to form the blue image, and although the part that the dots of both colors are overlapped becomes blue, the part which is not overlapped each other does not become blue, so that the deviation in color matching (irregular color) that each independent color tone appears is caused. When this occurs partially, it does not become easy to be seen, but when this phenomenon occurs in the direction of scanning continuously, a band-shaped deviation in color matching with a certain specific width is caused, so that the image becomes unequal. In addition, in a region adjacent the image region in the case of in the regions of the same color, when there is no offset in the depositing positions of the dots, a uniform impression and color development differ between the image regions adjacent each other, so that the image that there is a sense of incongruity as the image is formed. Moreover, though this deviation in color matching does not become easy to be seen in the case of an ordinary paper, it becomes easy to be seen, when a favorable printing medium in color development such as a coat paper is used.\nMoreover, in the case where a different color is printed on adjoining the pixel, when there is the offset in the depositing positions of the dot, the clearance, that is, the region which is not covered by the ink on the part have caused and, the ground of the printing medium can be seen. This phenomenon frequently is called xe2x80x9cwhite clearancesxe2x80x9d, since the case of a white ground is frequent in the printing medium generally. This phenomenon is easy to be seen in the image high in contrast, and when a black image is formed as a colored back ground, the white clearance which no ink is deposited between a black and coloration, since a contrast between white and black is high, can be easy to be seen more clearly.\nIt is effective to perform the above-mentioned dot alignment in order to suppressed occurring of the problems as mentioned above. However, the complicatedness that the user should observe the results which the depositing registration conditions are varied by the eyes to select the optimized the depositing registration condition to perform entering operations is accompanied, and moreover, since fundamentally, a judgment for obtaining the optimum printing position by observing through eyes is enforced on the user, the establishment which is not optimized can be set. Therefore, it is especially unfavorable to the user who is not accustomed to operation.\nMoreover, the user is enforced to expense in time and effort at least two times since the user should printing the image to perform the depositing registration and in addition, to perform conditional establishment after observing to perform judgments required, whereby upon realizing the apparatus or a system excellent in operability, it is not only desirable but also is disadvantageous from the viewpoint of a time-consumption.\nNamely, it has been desired strongly that the apparatus or system capable of printing the image at a high speed and of the high-quality image without occurring the problem on the image formation as above-mentioned and the problem on the operability is realized at a low cost by designing to be able to register the depositing position without using a feedback controlling means such as an encoder by an opened loop.\nTherefore, the object of the invention is to realize a dot alignment method which is excellent in operational performance and the low cost.\nMoreover, the invention, without fundamentally enforcing the user the judgment and the adjustment, is designed to detect the optical characteristics of the printed image to derive the adjustment condition of the optimum dot alignment from the detected results and to set the adjustment condition automatically, thereby to improve the adjustment accuracy thereof.\nIn the first aspect of the present invention, there is provided a printing registration method for performing a processing for performing a printing registration in a first printing and a second printing with respect to a printing apparatus for performing printing an image by the first printing and the second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, the method comprising the steps of:\nforming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeasuring optical characteristics for each of the plurality of patterns formed by the pattern forming step;\nacquiring adjustment value of dot forming position conditions between the first printing and the second printing on the basis of optical characteristics of respective of the plurality of patterns measured by the measuring step; and\nacquiring the adjustment value used for performing printing by executing plural times the pattern forming step, the measuring step and the adjustment value acquiring step for every different dot position registration accuracy.\nIn a second aspect of the present invention, there is provided a printing registration method for performing a processing for performing a printing registration in a first printing and a second printing with respect to a printing apparatus for performing printing an image by the first printing and the second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, the method comprising the steps of:\nforming patterns used for the processing by the first printing and the second printing;\nacquiring adjustment value of dot forming position conditions on the basis of the patterns; and\nperforming the step with every different dots position registration accuracy and acquiring an adjustment value used for performing the printing.\nIn a third aspect of the present invention, there is provided a printing registration method for performing a processing for performing a printing registration in a first printing and a second printing with respect to a printing apparatus for performing printing image by the first printing and the second printing with predetermined conditions of a dot forming on a printing medium by using a printing head, the method comprising the steps of:\nforming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and the second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeasuring a plurality of data by measuring plural times optical characteristics with positions made different with respect to each of the plurality of patterns formed by the pattern forming step;\nderiving optical characteristics of respective patterns from the plurality of data obtained per each of the plurality of patterns and acquiring adjustment value of dots forming position conditions between the first printing and the second printing on the basis of the optical characteristics; and\nmaking judgment per each of the plurality of patterns that the patterns are formed in a state not suitable for acquiring the adjustment value from the plurality of data obtained per each of the plurality of patterns and eliminating patterns formed in a state not suitable for acquiring the adjustment value from a processing of the value adjustment step.\nIn a fourth aspect of the present invention, there is provided a printing apparatus for performing printing an image by a first printing and a second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, comprising:\nmeans for forming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeans for measuring optical characteristics for each of the plurality of patterns formed by the pattern forming means;\nmeans for acquiring adjustment value of dot forming position conditions between the first printing and the second printing on the basis of optical characteristics of respective of the plurality of patterns measured by the measuring means; and\nmeans for acquiring the adjustment value used for performing printing by executing plural times of operations of the pattern forming means, the measuring means and the adjustment value acquiring means for every different dot position registration accuracy.\nIn a fifth aspect of the present invention, there is provided a printing apparatus for performing printing an image by a first printing and a second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, comprising:\nmeans for forming patterns used for processing of printing registration by the first printing and the second printing and for acquiring adjustment value of dot forming position conditions on the basis of the patterns;\nand wherein operations by the means are performed with every different dots position registration accuracy so that an adjustment value used for performing the printing is acquired.\nIn a sixth aspect of the present invention, there is provided a printing apparatus for performing printing image by the first printing and a second printing with predetermined conditions of a dot forming on a printing medium by using a printing head, comprising:\nmeans for forming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and the second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeans for measuring a plurality of data by measuring plural times optical characteristics with positions made different with respect to each of the plurality of patterns formed by the pattern forming means;\nmeans for deriving optical characteristics of respective patterns from the plurality of data obtained per each of the plurality of patterns and acquiring adjustment value of dots forming position conditions between the first printing and the second printing on the basis of the optical characteristics; and\nmeans for making judgment per each of the plurality of patterns that the patterns are formed in a state not suitable for acquiring the adjustment value from the plurality of data obtained per each of the plurality of patterns and eliminating patterns formed in a state not suitable for acquiring the adjustment value from a processing of the value adjustment means.\nIn a seventh aspect of the present invention, there is provided a printing system provided with a printing apparatus for performing printing an image by a first printing and a second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, and a host apparatus for supplying an image data to the printing apparatus, comprising:\nmeans for forming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeans for measuring optical characteristics for each of the plurality of patterns formed by the pattern forming means;\nmeans for acquiring adjustment value of dot forming position conditions between the first printing and the second printing on the basis of optical characteristics of respective of the plurality of patterns measured by the measuring means; and\nmeans for acquiring the adjustment value used for performing printing by executing plural times of operations of the pattern forming means, the measuring means and the adjustment value acquiring means for every different dot position registration accuracy.\nIn a eighth aspect of the present invention, there is provided a printing system provided with a printing apparatus for performing printing an image by a first printing and a second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, and a host apparatus for supplying an image data to the printing apparatus, comprising:\nmeans for forming patterns used for a processing of printing registration by the first printing and the second printing and for acquiring adjustment value of dot forming position conditions on the basis of the patterns; and\nwherein operations by the means are performed with every different dots position registration accuracy so that an adjustment value used for performing the printing is acquired.\nIn a ninth aspect of the present invention, there is provided a printing system provided with a printing apparatus for performing printing image by a first printing and a second printing with predetermined conditions of a dot forming on a printing medium by using a printing head, and a host apparatus for supplying an image data to the printing apparatus, comprising:\nmeans for forming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and the second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeans for measuring a plurality of data by measuring plural times optical characteristics with positions made different with respect to each of the plurality of patterns formed by the pattern forming means;\nmeans for deriving optical characteristics of respective patterns from the plurality of data obtained per each of the plurality of patterns and acquiring adjustment value of dots forming position conditions between the first printing and the second printing on the basis of the optical characteristics; and\nmeans for making judgment per each of the plurality of patterns that the patterns are formed in a state not suitable for acquiring the adjustment value from the plurality of data obtained per each of the plurality of patterns and eliminating patterns formed in a state not suitable for acquiring the adjustment value from a processing of the value adjustment means.\nIn a tenth aspect of the present invention, there is provided a storage medium which is connected to an information processing apparatus and a program stored in which is readable by the information processing apparatus, the program being for making a printing system to perform a printing registration method for performing a processing for performing a printing registration in a first printing and a second printing with respect to a printing apparatus for performing printing an image by the first printing and the second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, the method comprising the steps of:\nforming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeasuring optical characteristics for each of the plurality of patterns formed by the pattern forming step;\nacquiring adjustment value of dot forming position conditions between the first printing and the second printing on the basis of optical characteristics of respective of the plurality of patterns measured by the measuring step; and\nacquiring the adjustment value used for performing printing by executing plural times the pattern forming step, the measuring step and the adjustment value acquiring step for every different dot position registration accuracy.\nIn a eleventh aspect of the present invention, there is provided a storage medium which is connected to an information processing apparatus and a program stored in which is readable by the information processing apparatus, the program being for making a printing system to perform a printing registration method for performing a processing for performing a printing registration in a first printing and a second printing with respect to a printing apparatus for performing printing an image by the first printing and the second printing with predetermined conditions of a dot forming position on a printing medium by using a printing head, the method comprising the steps of:\nforming patterns used for the processing by the first printing and the second printing;\nacquiring adjustment value of dot forming position conditions on the basis of the patterns; and\nperforming the step with every different dots position registration accuracy and acquiring an adjustment value used for performing the printing.\nIn a twelfth aspect of the present invention, there is provided a storage medium which is connected to an information processing apparatus and a program stored in which is readable by the information processing apparatus, the program being for making a printing system to perform a printing registration method for performing a processing for performing a printing registration in a first printing and a second printing with respect to a printing apparatus for performing printing image by the first printing and the second printing with predetermined conditions of a dot forming on a printing medium by using a printing head, the method comprising the steps of:\nforming a plurality of patterns respectively having optical characteristics corresponding to a plurality of shifting amounts, the plurality of patterns being patterns formed by the first printing and the second printing by using the print head, and the plurality of patterns being formed respectively corresponding to the plurality of shifting amounts of relative printing positions of the first printing and the second printing;\nmeasuring a plurality of data by measuring plural times optical characteristics with positions made different with respect to each of the plurality of patterns formed by the pattern forming step;\nderiving optical characteristics of respective patterns from the plurality of data obtained per each of the plurality of patterns and acquiring adjustment value of dots forming position conditions between the first printing and the second printing on the basis of the optical characteristics; and\nmaking judgment per each of the plurality of patterns that the patterns are formed in a state not suitable for acquiring the adjustment value from the plurality of data obtained per each of the plurality of patterns and eliminating patterns formed in a state not suitable for acquiring the adjustment value from a processing of the value adjustment step.\nOptical characteristics (characteristics of changes in density) with respect to the dot formative positions condition is changed based on the relation of pixel density and a dot diameter, depending upon a formation positions of the dot greatly, whereby from the characteristics the relative dot-formation position can be obtained.\nThe condition that the dots which are adjacent are in contact with each other is largest in planar dimension, as it approaches from a connecting condition, the planer dimension is decreased in accordance with a change of the formation position. In other words, the density is changed in accordance with the formation position. Moreover, from the relation of the pixel density and a dot diameter, in order to make the area factor to 100%, the dot has a diameter of size of {square root over (2)} times of one pixel, and under the condition that the formation position is registered the overlapped parts exist inescapably in the dots which are adjoined each other, and at that condition, the density becomes maximum. In contrast with this, the formation position is deviated, whereby when the condition that the area factor does not become 100%, that is, the condition which a clearance can be formed is achieved, the density is decreased.\nTherefore, the condition that the formation position are registered is the region where the density is changed greatly in the formation position of the dot. By varying the position registration condition of the formation position of the dot with respect to the dot as the reference the density is made change. Using the change in density, approximate characteristics are obtained, so that the adjusting position where the depositing position of the dot have just registered.\nFurther, it is possible to obtain an adjustment range over a maximum offset predicated by accuracies of a printing apparatus and a printing head, and to execute a processing including a coarse adjustment to a fine adjustment in a series of algorithm. That is, by performing an adjustment method according invention, it is possible to adjust a printing apparatus where dot forming position are offset largely, to the state in which optimum dot forming positions are obtained, without user\"\"s decision.\nThis adjustment method is adapted to the strict adjustment of the depositing position, and a dot alignment (a printing registration) with high accuracy can be realized, since the slight offset of the formation position appears sensitively on the change in density. Further, the method is applicable to a printing apparatus having a large offset of deposition position accuracy, and realize a dot alignment with a wide adjustment range.\nIncidentally, hereafter, the word xe2x80x9cprintxe2x80x9d (hereinafter, referred to as xe2x80x9crecordxe2x80x9d also) represents not only forming of significant information, such as characters, graphic image or the like but also represent to form image, patterns and the like on the printing medium irrespective whether it is significant or not and whether the formed image elicited to be visually perceptible or not, in broad sense, and further includes the case where the medium is processed.\nHere, the wording xe2x80x9cprinting mediumxe2x80x9d represents not only paper to typically used in the printing apparatus but also cloth, plastic film, metal plate and the like and any substance which can accept the ink in broad sense.\nFurthermore, the wording xe2x80x9cinkxe2x80x9d has to be understood in broad sense similarly to the definition of xe2x80x9cprintxe2x80x9d and should include any liquid to be used for formation of image patterns and the like or for processing of the printing medium."}
{"text": "1. Field of the Invention\nThe invention relates to a tool saver for use in association with oil or gas wells or the like. More specifically, the invention relates to a tool saver which prevents tools, when withdrawn from a borehole of a well, from falling back into the well in the event of a break in the wire line which lowers and raises the tools.\n2. Description of Prior Art\nAs is known in the art, wire line tools for wells, such as, for example, sand line swabbing equipment, sinker bars, etc., are lowered into and raised out of wells through the well bores of the wells by wire lines. There is a constant danger that, after the tool is raised out of the well, the wire line will break and the tool will fall uncontrollably back into the well. The tool will then have to be retrieved at great cost and inconvenience.\nVarious devices are known for gripping tools or pipes as they are lowered into wells. Examples of such devices are taught in, for example, U.S. Pat. Nos. 1,422,289, Moody, Jul. 11, 1922, 4,076,337, Childress, Feb. 28, 1978, 4,096,608, Lagerstedt, Jun. 27, 1978 and 4,898,238, Grantom, Feb. 6, 1990.\nThe '289 patent teaches a device for gripping casings or sections of pipes as they are being raised or lowered into a well. Specifically, the device will hold one section while another section is being coupled or uncoupled.\nThe drill steel holder of the Childress patent, U.S. Pat. No. 4,076,337, comprises levers 26 and 28 (see FIGS. 3 and 4) having outboard ends connected to a RAM 40. The RAM 40 can draw the outboard ends towards each other or force them apart from each other.\nThe inboard ends of the levers 26 and 28 comprise rollers 50. When the inboard ends of the levers are drawn towards each other, as shown in FIG. 3, the rollers will embrace a tool, for example a steel drill 68, to prevent the vertical movement of the steel drill.\nJaws 24 and 25 of U.S. Pat. No. 4,096,608, (having inserts 28 and 29) are forced towards each other by a spring 49 to capture a drill pipe 19 as shown in FIG. 2 of the drawings. As seen in FIG. 1, this will prevent the drill pipe 19 from moving in a vertical direction.\nIn U.S. Pat. No. 4,898,238, Grantom, a stationary jaw 200 acts together with a gripping, or movable, jaw 300 to capture a pipe or the like which is inserted into a well. As seen in FIG. 1, the movable jaw 300 is moved by the action of threaded shaft 320 in threaded block 310.\nAs is quite clear, none of the references address themselves to the problem as above described."}
{"text": "This invention relates generally to apparatus and methods for making induced polarization and resistivity measurements in a borehole.\nIn the study of geological formations it is known that induced polarization and resistivity measurements can be made by means of surface-located equipment. See, for example, U.S. Pat. No. 4,295,056. Through the use of such surface equipment, induced polarization measurements which may include phase angle anomalies can be obtained. The phase angle anomalies can be created by a shallow subsurface polarizable material which is of significance to the oil and gas industry if its composition can be geochemically tested and found to contain hydrocarbons which have seeped from a deeper reservoir. However, before the material can be tested, its location must be pinpointed. Therefore, there is the need for an apparatus and a method which utilize electrodes in a borehole to make downhole induced polarization and resistivity measurements in the vicinity of the formations which are creating the anomalies so their locations can be accurately determined. Once located, the formations can be geochemically tested to determine if hydrocarbons are present (see, for example, U.S. Pat. No. 4,360,359). From such information one can tell whether hydrocarbons are seeping from a reservoir, for example. The need for an apparatus and a method for making such downhole measurements has been recognized. See, for example, Complex Formation Resistivity--The Forgotten Half of the Resistivity Log by D. D. Snyder, R. H. Merkel, and J. T. Williams and Induced Polarization Measurements In and Around Boreholes by D. D. Snyder and R. H. Merkel.\nIt has been observed that some of these anomalies arise from shallow or thin formations, thereby making it desirable to utilize an apparatus which can make measurements of sufficient resolution for such thin formations. This desirability indicates that a short array spacing of the electrodes be used, which short array spacing corresponds to the size of the thin layer; however, for such a short array spacing of the type needed to make the induced polarization and resistivity measurements, the measurements made thereby would be adversely affected by the electrical properties of the drilling mud or other fluid surrounding the electrodes at their measurement locations in the borehole. That is, a short array spacing would cause the apparatus to read the electrical properties of the drilling mud or other fluid rather than the electrical properties of the desired formation.\nThe borehole fluid can have a substantial effect on the phase readings made by such an apparatus because the borehole fluid may be a drilling mud typically containing highly polarizable montmorillonite clay, for example. Additionally, it is difficult, if not virtually impossible, to determine ahead of time exactly what the effect of the borehole fluid will be. Still further, a drilling mud tends to be electrochemically unstable so that it is difficult to determine at any given time what the electrical properties of the mud are. Furthermore, the depth of invasion of the mud or other borehole fluid into the rock formation is highly variable. Therefore, there is the need for the apparatus and method to be constructed and performed so that several suitable measurements with variable length arrays can be made and compared to derive a reasonable estimate of the electrical properties of thin formations. This reasonable estimate can be further refined by suitable modeling of the formation surrounding the borehole.\nIt is also desirable that the apparatus and method which meet the foregoing needs be usable with commercially available interconnection equipment presently found in the oil and gas industry. This commercially available equipment includes electrical logging cable comprising parallel electrical conductors. Such commercially available cable is susceptible to significant electromagnetic coupling between the various parallel conductors. This electromagnetic coupling adversely affects the induced polarization measurements.\nOne way to avoid the electromagnetic coupling is to use separately twisted pairs of wires; however, such twisted pairs of wires are not standard equipment in the oil and gas industry and thus their use would require special orders with possibly long lead times and high costs. Additionally, separately twisted wires would make for a thick cable which would reduce the amount of cable which could be retained on available winches.\nSuch a problem of electromagnetic coupling could also be overcome by placing certain equipment in the implement which is to be used in the borehole and telemetering the readings to the surface. A shortcoming of placing additional equipment in the borehole is that it requires packaging the equipment for use in the extreme pressures and temperatures found in the borehole. Furthermore, having the equipment downhole would limit the ability to control it, which ability would be present if the equipment were to be located at the surface. Telemetering signals to the surface would limit the dynamic range and accuracy of the measurements made in the borehole. Therefore, there is the need for the apparatus and method to be constructed and performed so that such adverse electromagnetic coupling is eliminated without preventing the apparatus and method from being used with or implemented by commercially available equipment.\nIt is known that such induced polarization and resistivity measurements as are referred to hereinabove are made by introducing a current into the formation and then detecting a voltage produced in the formation by the current. It has been discovered that it is important to control the density of the current because if the current density is too large, the induced polarization response will be non-linear with respect to the current and thus will not be representative of the response that a surface induced polarization survey would obtain. Therefore, the apparatus and method which meet the aforementioned needs should also meet the needs of maintaining an accurately controlled current output and utilizing a low-noise receiver which can accommodate very small transmitted currents.\nTo insure that accurate measurements are made, the apparatus should be constructed so that electrical leakage in the portion of the apparatus which is to be used in the borehole is reduced. It has been found that the phase measurement which is indicative of the induced polarization response is much more sensitive to minor amounts of electrical leakage than are conventional magnitude-only resistivity measurements. To prevent such electrical leakage, there is the need for the apparatus to be constructed so that there is substantially a zero pressure differential between the interior of the downhole implement and the borehole pressure to thereby preclude the flow of borehole fluid into the downhole implement.\nTo further insure the accuracy, the desired method is to be performed discretely whereby discrete, rather than continuous, measurements of the induced polarization and resistivity are made. Discrete measurements preclude the introduction of noise and instability which can occur by means of the movement of the downhole implement through the borehole fluid as measurements are taken. Additionally, the preferred current used to obtain such induced polarization and resistivity measurements is at a very low frequency so that if a continuous logging run were made, it could cover a substantial distance before one cycle of the current were transmitted.\nMore generally, there is the need for an apparatus and a method for obtaining measurements which can be used not only to qualitatively look for anomalies, but also to quantitatively understand surface-obtained induced polarization data, other surface geological and geophysical data and downhole electrical properties."}
{"text": "Molecular biotechnology is a discipline that is based on the ability of researchers to transfer specific units of genetic information from one organism to another with the goal of producing commercially relevant amounts of useful products. One of the goals of this cloning process is to achieve maximum expression of the cloned gene. Recombinant production of a product encoded by a gene is accomplished by constructing expression vectors suitable for use in a host cell in which the nucleic acid coding for a desired product is placed under the expression control of a promoter. The expression vector is introduced into a host cell by various techniques, such as transformation, and production of the desired product is then achieved by culturing the transformed host cell under suitable conditions necessary for the functioning of the promoter included in the expression vector."}
{"text": "Generator rotors are provided with radial slots about the periphery thereof, for receiving field or rotor windings of coils made up of a number of turns in a radially stacked arrangement, each turn or winding separated by insulation. The windings are retained in each slot by one or more full-length wedges, with one or more creepage blocks interposed between the wedges and the windings. The function of the wedge block is to contain the contents of the rotor slot against outward radial forces arising from rotation of the rotor. The creepage block serves as an insulating barrier between the wedge and the copper windings. The creepage block is also formed with apertures which permit cooling gas to flow from radial passages in the copper windings to radial passages in the wedge.\nGenerators currently available from the assignee are placed in three major design classifications based on the cooling medium used: air, hydrogen or liquid (water). All hydrogen and water cooled generators use direct conductor cooling of the rotor windings for heat removal. Smaller two-pole and all four-pole generators use the radial flow design.\nIn such radial flow systems, the field windings imbedded in the rotor slots are ventilated by a radial cooling circuit consisting of an axial sub-slot and a series of radial cooling passages. To maximize cooling gas flow in the radial direction and electrical conductance in the axial direction of the winding, the axially spaced radial flow cooling ducts are elongated in the axial direction. To minimize stresses in the rotor wedges, the last leg of each radial cooling path is a round hole. The problem to be solved is a cost-effective optimization of the transition via the creepage block from elongated ducts in the copper windings to the circular holes in the wedge without incurring undesirable pressure losses that would undermine the thermal performance of the ventilation circuit.\nHistorically, the creepage block radially between the wedge and the coil windings has been made by introducing a multitude of geometric features for each radial passage. In one example, cylindrical holes are drilled concentric with each wedge hole. Circular arches are then machined on the radial inboard surface of each creepage block to match the elongated ducts in the windings. Finally, a shallow axial groove is machined through the radially outboard surface of each wedge to allow for small variations in alignment between the creepage block and the wedge block. This conventional arrangement is illustrated in FIG. 1 (described in detail further below)."}
{"text": "Embodiments of the present invention relate to a non-volatile memory device and method for controlling the same, and more specifically, to a technology for increasing the operation reliability of a non-volatile memory device in such a manner that the sensing and transmitting of data can be stably achieved during a read operation in which data stored in a unit cell is output externally.\nSemiconductor memory devices have been rapidly developed to temporarily or permanently store data therein. The semiconductor memory devices have been widely used in a variety of electronic appliances, electronic devices, and personal portable devices. General semiconductor memory devices can freely read and/or write data, and can also easily update old data with new data.\nThe semiconductor memory device has been increasingly developed such that it can store much more data therein, can be operated with a smaller amount of power, and can have an increased or expedited operating speed. Although a NOR flash memory device or a NAND flash memory device has been widely used as non-volatile memory, a conventional flash memory device has a disadvantage in that it has a slow operating speed. In order to overcome the above-mentioned disadvantage, a non-volatile solid status memory device, generally called a magnetic random access memory (RAM) device or a magnetic memory device, has been recently proposed and developed.\nA non-volatile memory device formed of a magnetic material has not only characteristics of a dynamic RAM (DRAM) device which has a high data processing speed, a high integration degree, and low power consumption, but also characteristics of a flash memory device capable of preserving data even when not powered. Unlike a DRAM device for storing information according to an electrical charging scheme, the non-volatile memory device formed of the magnetic material stores information or data using a magneto resistive effect, and may be implemented with a Giant Magneto Resistive (GMR) element or a Tunneling Magneto Resistance (TMR) element.\nFIG. 1 is a circuit diagram illustrating a read operation of a general non-volatile memory device.\nReferring to FIG. 1, the general non-volatile memory device includes a unit cell 110, a bit line 120, a source line 130, a sense-amplifier (also called a sense-amp) 140, a first switching unit 150, a second switching unit 160, a third switching unit 170, and a fourth switching unit 180.\nIn this case, the unit cell 100 includes a Magnetic Tunneling Junction (MTJ) element 112 and a cell transistor 114. The MTJ element 112 has a resistive MTJ (RMTJ) value, and the cell transistor 114 is controlled by a voltage of a word line WL.\nThe bit line 120 is coupled to the MTJ element 112, such that it has a bit line resistance (RBL) value. The source line 130 is coupled between one end of the cell transistor 114 and the second switching unit 160, such that it has a source line resistance (RSL) value.\nThe sense-amp 140 includes a current source IREF, a switching element, and an amplifier, such that it senses and amplifies data stored in the unit cell 110. In this case, the current source IREF generates a bit line current IBIT in response to a core voltage VCORE. The amplifier compares a voltage generated by the current source IREF with a reference voltage VREF, such that it outputs an output voltage VOUT. The switching element is controlled by a voltage suppression signal (also called a voltage clamping control signal) VCMP, such that it restricts a level of a bit line current IBIT generated by the current source IREF.\nGenerally, a memory is configured in the form of an array which is composed of a plurality of cells. Therefore, on a current path to obtain data stored in the specific unit cell 110, there are parasitic resistances of common signal lines 120 and 130 and the first switching unit 150, the second switching unit 160, the third switching unit 170, and the fourth switching unit 180 for selecting the specific unit cell 110.\nThe first switching unit 150 controls a connection between the bit line 120 and the third switching unit 170 when selecting the unit cell 110 in response to a selection signal SEL serving as an output signal of a decoder. In addition, the second switching unit 160 controls a connection between the source line 130 and the fourth switching unit 180 when selecting the unit cell 110 in response to the selection signal SEL. The third and fourth switching units 170 and 180 control a read operation in response to a read control signal RDE from an external part.\nIn order to show how the current IBIT flows in a path including the unit cell 110 when the non-volatile memory device performs the read operation, FIG. 1 illustrates resistance values belonging to all constituent elements.\nIn the non-volatile memory device such as a magnetic memory, data is stored as a resistance value in the MTJ element 112. Generally, the MTJ element 112 has a three-layered structure in which one insulation layer is located between two ferromagnetic layers. In the non-volatile memory device, a magnetization direction of the ferromagnetic layer is determined depending on a direction of a current applied to the MTJ element 112, such that a resistance value is determined depending on the current direction.\nIn order to obtain the data stored in the magnetic memory, a given voltage is applied to both ends of the MTJ element 112, and the current flowing in the MTJ element 112 is measured, so that it is determined whether the MTJ element 112 is in a high resistance (RH) state or a low resistance (RL) state.\nFor example, digital data of ‘0’ or ‘1’ may be determined in response to a resistance value of the MTJ element 112. Such digital data ‘0’ and ‘1’ may respectively represent a case in which the resistance value of the MTJ element 112 is higher than a reference resistance value and the other case in which the resistance value of the MTJ element 112 is equal to or lower than the reference resistance value. This allows a user to recognize which data is stored in the MTJ element 112.\nIn other words, the first case in which the MTJ element 112 has a high resistance RMTJ value means that two magnetic substances constructing the MTJ element 112 have opposite magnetization directions. On the other hand, the second case in which the MTJ element 112 has a low resistance RMTJ value means that the two magnetic substances have parallel magnetization directions.\nIn this case, although the resistance RMTJ value of the MTJ element 112 contained in the unit cell 110 is changed according to data, resistance values of constituent elements other than the resistor RMTJ are not changed depending on data. That is, resistance values of the first to fourth switching units 150-180 for selecting the unit cell 110 or performing the read operation, a resistance value of the bit line 120, and a resistance value of the source line 130 have fixed values that are not changed depending on data.\nIn the case of performing the read operation, a sensing voltage VRD for detecting data is applied to one end of the unit cell 110 and a ground voltage is applied to the other end of the unit cell 110, such that the current IBIT flowing in the unit cell 110 is measured. In this case, the sense-amp 140 compares the measured current IBIT with a reference value VREF, such that it determines whether a resistance value of the MTJ element 112 is higher or lower than the reference value VREF.\nIn a general non-volatile memory device, the sense-amp 140 is not individually coupled to one unit cell 110, and is coupled to a cell array including a plurality of unit cells.\nIn addition, during the read operation, the sensing voltage VRD passes through not only one unit cell 110 but also a plurality of constituent elements having unique resistance values, such that the current IBIT flowing in the unit cell 110 is very small when the read operation is performed.\nThe small current IBIT means that the intensity of a signal to be output in response to data is very low. A TMR value is used to define the signal intensity. The TMR value is defined as a specific value obtained when a difference between resistance values of two states of the MTJ element 112 is divided by a smaller one of the two resistance values. The TMR value may be represented by “TMR=(RH−RL)/RL*100[%]”.\nHowever, assuming that one MTJ element 112 in the non-volatile memory device has a TMR value of about 100%, when measuring the TMR value of the MTJ element 112 along a read path of a cell array contained in an actual non-volatile memory device, the TMR value of about 100% is reduced to a TMR value of about 30% due to parasitic resistance of other constituent elements.\nIn response to the voltage clamping control signal VCMP, the sense-amp 140 may apply the core voltage VCORE to the read path so as to restrict the sensing voltage VRD to a predetermined level or less, such that it substantially prevents an excessive current from being applied to the MTJ element 112 and thus physical properties of the MTJ element 112 are not broken or deteriorated.\nFIG. 2 is a graph illustrating the current IBIT flowing in the read operation of the non-volatile memory device shown in FIG. 1.\nReferring to FIG. 2, the sensing voltage VRD is controlled to be a predetermined level or less in response to the voltage clamping control signal VCMP. And, when the read control signal RDE, the selection signal SEL for selecting the unit cell 110, and the word line WL are activated, the current IBIT flows along the read path in response to data stored in the unit cell 110.\nThus, the current is changed depending on information stored in the MTJ element 112. That is, if the MTJ element 112 has high resistance, the current IBIT flowing in response to data is changed to a low current IH smaller than a reference current IREF. Otherwise, if the MTJ element 112 has low resistance, the current IBIT flowing in response to data is changed to a high current IL greater than the reference current IREF. In this case, the low current IH and the high current IL are not always set to a fixed value, and may change according to an operation environment.\nIn FIG. 2, the horizontal axis represents data current IBIT, and a vertical axis represents a relative frequency P(I) of a cell having the data current IBIT.\nAs shown in FIG. 2, the sense-amp 140 determines whether or not the current IBIT is greater than the reference current IREF, such that it recognizes data stored in the unit cell 110. In this case, in order to substantially prevent an error from being generated in the read operation, it is necessary for the reference current IREF to maintain the range from the low current IH to the high current IL. In order to maximize an operation margin of the sense-amp 140, it is necessary for the reference current IREF to have an intermediate value between the low current IH and the high current IL.\nIn the case of using an absolute reference current scheme shown in FIG. 2, the reference current IREF maintains a predetermined value, such that a circuit for generating the reference current IREF can be simplified in structure. In this case, in order to maximize the operation margin of the sense-amp 140, the cell array needs to be designed in such a manner that the relative frequency P(I) of the current IBIT is maximized.\nHowever, it should be noted that the P(I) distribution is not always maintained in an invariable shape. The data current distribution (i.e., variation of a position of P(I) and shape variation) may be generated by a variety of reasons. A representative one of such reasons is variation in resistance existing on the read path illustrated in FIG. 1.\nSpecifically, the relative frequency P(I) may be affected by where the specific unit cell 110 selected in the read operation is in the cell array. That is, the bit line resistor RBL, the source line resistor RSL, resistance of transistors in the first to fourth switching units 150˜180 serving as various kinds of switches, and resistance generated when the sensing voltage VRD is provided may affect the selected unit cell 110. In addition, in the process for manufacturing the non-volatile memory device, resistance variation caused by process variations of the MTJ element 112 may also encounter the variation of the P(I) position and shape.\nIn the case where data is sensed and amplified using the reference current IREF having an invariable value on the condition that the P(I) position and shape are changed due to the aforementioned reasons, the operation margin of the sense-amp 140 may be reduced. If worst comes to worst, a faulty operation in which data is incorrectly sensed may be generated.\nFIG. 3 is a graph illustrating the current IBIT in the read operation of the non-volatile memory device so as to solve the problems shown in FIG. 2.\nIn order to increase the operation margin of the sense-amp 140 capable of sensing and amplifying the current IBIT flowing in the unit cell 110 during the read operation, a relative reference current scheme is used, instead of the absolute reference current scheme in which the reference current IREF has a predetermined value. The relative reference current scheme may also change the value of the reference current IREF in response to the variation of the relative frequency P(I) of the current IBIT.\nThe relative reference current scheme can maintain the sensing margin, such that it can readily read low-intensity information of the unit cell 110. However, in the relative reference current scheme, both the area of a circuit generating the reference current IREF and power consumption are increased so as to change the reference current IREF."}
{"text": "Lottery tickets are well known and widely sold and typically comprised of a sheet material of paper or card stock on which is printed lottery information and various indicia for the playing of one or more games. Many such games are instant win type games where the player can play the game or games by carrying out various functions, for example, opening pull tabs on a break-open ticket. Such tickets are also known variously as pull-tabs, pickle cards, jar tickets, hard cards and charitable gaming tickets.\nBreak-open tickets are typically constructed by printing a sheet containing multiple combinations of lottery indicia thereon. This sheet is then cut into sections. The opposite side of this sheet includes prize categories, serial number etc. A second sheet which contains removable tabs or break-open windows is then mounted onto the card stock such that the break-open windows overlap the lottery indicia. However, break-opens prepared in this manner have the disadvantages that they are cumbersome to prepare, have a limited set of winning symbols and require manually checking at the time of redemption.\nU.S. Pat. No. 4,174,857 (Koza) and U.S. Pat. No. 5,253,899 (Greenwood) each disclose a combined ticket construction in which a second layer is applied or folded over the basic ticket to improve security. U.S. Pat. No. 5,562,284 discloses a break-open ticket with a second level of break-open. U.S. Pat. No. 5,217,258 (Sanderson) discloses a two sided break-open ticket where the windows are arranged back to back opening onto a single center sheet carrying the game data.\nAs shown in published PCT International application No. WO 98/57715 published 23 Dec. 1998 of the present applicants there is provided a pouched lottery ticket which is supplied in a fan folded continuous strip arrangement for separating each ticket from the next at a dispensing station. The disclosure of this publication may be referred to for further information.\nIn U.S. Pat. No. 6,234,477 assigned to the present assignee and issued May 22, 2001 is disclosed a construction of lottery ticket which utilizes variable image printing techniques in combination with lamination of two substrate sheets together to form what is initially a common pouch construction with the game indicia on the inside surface where the sheets can be separated, each from the other, to form in effect two separate tickets."}
{"text": "“Alkylation” generally refers to the reaction of a hydrocarbon, such as an aromatic or a saturated hydrocarbon, with an olefin. For example, in one type of reaction of particular interest, a branched saturated hydrocarbon, such as isobutane, may undergo alkylation with an olefin containing 2-6 carbon atoms, such as 2-butene, to produce an alkylate that has a higher octane number and which boils in the gasoline range. Processes directed to the alkylation of paraffins with olefins produce branched hydrocarbon molecules for gasoline components, such as isomers of octane, e.g. trimethylpentanes (“TMPs”), which have high octane numbers. A gasoline with a high octane number, often expressed as research octane number (“RON”), can reduce engine knock, which lessens the need to add environmentally harmful anti-knock compounds such as tetraethyllead. A second octane measurement, motor octane number (“MON”), also describes the anti-knock properties of gasoline. MON is measured when the test engine is run under heavy load (higher rpm), and the RON is measured at lower load (lower rpm).\nGasoline produced by the alkylation process is essentially free of contaminants, such as sulfur and nitrogen that may be present in gasoline obtained by other processes, such as cracking heavier petroleum fractions, e.g. vacuum gas oil and atmospheric residue. Sulfur oxides (“SOx”), a combustion product, are the primary cause of pollutants. In addition to direct SOx emissions, SOx can significantly lower the effectiveness of catalytic converters, thereby adversely impacting SOx, NOx, and CO emissions. SOx also form indirect particulates—a combination of water and SOx to form sulfurous and sulfuric acids. These indirect particulates normally exist in the 1-10 micron range, which are “inhalable particulates” that cause health problems, especially to people who suffer from asthma or emphysema. Also, unlike gasoline obtained by reforming naphtha or by cracking heavier petroleum fractions, alkylate contains few if any aromatics or olefins. Aromatics, especially benzene, are toxic, and olefins are reactive in photochemical reactions which cause ozone and smog.\nThe alkylation reaction is acid-catalyzed. Liquid acid catalysts, such as sulfuric acid or hydrofluoric acid, have been commonly used in alkylation processes. The use of liquid acid catalysts has several disadvantages. The liquid acids used are highly corrosive, requiring special quality, more expensive equipment. Because the presence of these acids in the resulting fuel is undesirable, any acid remaining in the alkylate must be removed. This process is complicated and expensive. In addition, the liquid acids, especially hydrofluoric acid, are dangerous if released into the environment.\nTo address these and other deficiencies of liquid acid catalysts, solid acid catalysts have been developed for use in alkylation processes. The solid catalysts typically employ a solid acid catalyst and a metal that provides a hydrogenation function. For example, U.S. Pat. No. 6,855,856 describes a catalyst comprising a solid acid, such as a zeolite, and a hydrogenation function. The solid acid described has a defined range for the ratio of the volume of catalyst pores to the specific length of the catalyst particles.\nA disadvantage of the prior solid acid catalysts is that the catalyst can become rapidly deactivated due to the formation of polyalkylates (e.g. C12+ product) which inhibits the alkylation reactions—somewhat like very soft coke. As soon as the catalyst forms a certain level of polyalkylates, the catalyst essentially stops the alkylation reactions. In a fixed bed reactor, an often preferred configuration, one can view the deactivation as occurring as a band-wise aging, with the deactivation zone moving as a band throughout the bed until most of the bed is inactive. This catalyst deactivation requires that the catalyst be periodically regenerated to ensure that the process produces a sufficient yield of the desired product. Regeneration of the catalyst typically requires that the alkylation process be stopped for a period of time. This reduces production and increases the cost of the alkylation process, especially by lowering the “onstream” factor of the process.\nA preferred method of regeneration of the catalyst is hydrogenation. The hydrogenation function is typically provided by a metal of Group VIII of the Periodic Table of the Elements, in particular the noble metals such as platinum (Pt) or palladium (Pd). Unlike the classical bifunctional (metal/acid) catalyst, the hydrogenation function plays little or no direct role in the alkylation reactions itself. Instead, it plays a critical role in the effective H2 reactivation (also called “regeneration” here) of the deactivated catalyst. The hydrogenation function is important in both the so-called low temperature (“low T”) and high temperature (“high T”) regenerations, described below.\nVarious attempts have been made to develop improved solid acid catalysts. For example, U.S. Publication No. 2004/0162454 describes an alkylation catalyst comprising nanocrystalline zeolite Y and a hydrogenation metal. The pore size of the nanocrystalline zeolite Y provides an alkylate with a higher RON/MON, as well as a longer run time for the catalyst. The nanocrystalline zeolite Y catalyst also includes a metal of Group VIII of the Periodic Table of the Elements, such as Pt or Pd, to provide a hydrogenation function.\nTo increase the efficiency and productivity of the alkylation process using solid acid catalysts, various methods have been developed to improve the process of regenerating solid acid catalysts. For example, U.S. Pat. No. 7,176,340 describes a continuous process for alkylation using a total of at least four catalyst containing reactors. However, use of multiple reactors increases the cost of the process; this cost increase may be offset, at least in part, by the efficiency increase of the overall process U.S. Pat. No. 5,986,158 describes an alkylation process in which the catalyst is subjected intermittently to a regeneration step by being contacted with a feed containing a saturated hydrocarbon and hydrogen, with the regeneration carried out at 90% or less of the active cycle of the catalyst. While these regeneration methods improve the overall efficiency of the alkylation process, the relatively large amounts of solid acid catalysts and associated noble metals required could be a problem that affects the commercial viability of the alkylation process.\nIt would be desirable to have a solid acid catalyst for the alkylation process that provided longer run times prior to deactivation. It would also be desirable to have a solid acid catalyst that that utilizes a metal for the hydrogenating function that provides equal or improved performance as compared to Pt or Pd, and that may be available at a lower cost. The present invention overcomes one or more of these and other drawbacks or disadvantages of prior solid acid catalysts used in alkylation processes."}
{"text": "Locating systems are becoming widespread with the onslaught of location-based services. Applications that depend on such services are growing, and the need for accurate measurement information is ever present and compelling."}
{"text": "Embodiments of this disclosure will describe devices and methods for conducting measurements to determine status and/or health of tissue. One area of the disclosure will address neuromuscular diseases, which are very common. In 2009, there were over 18 million outpatient physician encounters with patients diagnosed with a range of neuromuscular diseases, from ALS to myoneural disorders like myasthenia gravis. These episodes generated nearly $7B in physician charges and can only be expected to grow as the episode volume is expected to reach over 20 million by 2014 (+8.25%).\nAnother area our disclosure will address common medical complaints, such as, e.g., lower back and neck pain. These common medical complaints are sometimes the primary causes of disability, lost productivity, and medical costs. For example, 60-80% of adults experience at least one significant episode of back pain in their lifetime. In a single year, about 15% will have debilitating back pain, and a substantial proportion will seek medical attention. Lower back pain has been estimated as the fifth or higher leading cause of all medical visits, and the first or second leading cause for patients seeking evaluation and treatment of a condition. Stated another way, about 4-5% of all medical encounters are related to back pain. In two major health surveys (National Health Interview Survey and National Health and Nutrition Survey), from 2004-2008, 28-40% of the U.S. population experienced neck or back pain in a three-month period, with 14-21% experiencing neck pain. Of the total group, 26-33% experienced associated radicular pain in a limb. Remarkably, lumbosacral and cervical pain together caused 5% of all U.S. health care visits in the in 2006.\nCosts associated with care of individuals with low back and neck pain are huge. For example, in 2006, 44.4 million patients sought medical attention for low back pain, which was the chief complaint in 45.1 million encounters; an additional 13.2 million medical encounters were for neck pain. Annual cost of back pain in the US is $20-50 billion. Direct medical costs in 2002-2004 for spine problems were $193.9 billion, with $30.3 billion attributed specifically to spine pain. Indirect costs of lost wages are estimated at $14 billion annually, and in 2008, 385 million work days were lost due to back pain. The subgroup of patients with pain radiating to the limb, as occurs with disc herniation or spinal stenosis, had the highest numbers of bedridden and lost work days. Non-physician health care visits, e.g., physical therapy and other services, numbered 173.5 million from 2002-2004.2 Utilization is increasing, with ambulatory physician visits up 2.5% and non-physician visits up 10.2% from 1996-1998 to 2002-2004, with total increase in health care costs for spine of 24.5% (mean) and 48.9% (aggregate) in the same time period.\nAnother area our disclosure will address is the muscular health of older adults, which is the fastest growing segment of the population. Health care cost are exceptionally high for older adults with declining function, accounting for a disproportionate fraction of national health care expenditures. Muscle weakness is an independent risk factor for disability and mortality among older adults. Age associated loss of muscle mass, known as sarcopenia, is an important factor identified as relevant to mortality and disability.\nIn a recent and important editorial on sarcopenia and muscle function, Ferrucci et. al. (L. Ferrucci, R. de Cabo, N. D. Knuth, S. Studenski “Of Greek heroes, wiggling worms, mighty mice and old body builders”, J. Gerontol A Biol Sci Med Sci 2012; 67:13-6, [Ferrucci 2012] which is incorporated herein in its entirety by reference) advocated for a clinical approach in evaluating age associated muscle impairments, stating that after an initial mobility assessment to stratify patient risk for adverse outcomes (e.g., disability and mortality), muscle strength should be measured and a decision tree assessment used to evaluate muscle quality and function.\nBoth the general research community and the U.S. Food and Drug Administration (FDA) recognize the importance of improved biomarkers for neuromuscular disease research to assist with early diagnosis and track disease progression over time and response to therapy. Even more fundamentally, the concept of biomarker has expanded beyond its earlier definition that was restricted to molecular indices and now includes, but it not limited to, imaging and other methodologies. In fact, the FDA defines a biomarker as any objective test of disease status that cannot be influenced by the state of mind of the patient or examiner. The FDA recently developed biomarker definitions including diagnostic biomarkers for disease identification, response predictive biomarkers for assessing subgroups of individuals more likely to respond to a specific therapy, and prognostic biomarkers, for evaluating likelihood of disease onset or progression without any form of intervention. The FDA definitions provided herein are for discussions and references purposes only, and are not intended to limit any term contained herein. Two categories of biomarkers include response identification biomarkers (also called pharmacodynamic biomarkers) and biomarkers as surrogate endpoints in clinical trials. As a result, the FDA is revamping its approach to drug approval based on such surrogate endpoints. Previously, approval demanded evidence of change in a clinical outcome measure, such as, e.g., improved physical function or activity. In the future, however, it may be possible for a biomarker, which is established as a surrogate endpoint in a clinical trial, to obtain “qualification” status through the FDA as a surrogate endpoint, helping speed study and approval of effective therapies.\nOne technique for evaluating muscles is intramuscular electromyography (EMG.) EMG includes, but is not limited to, a technique for evaluating and recording the electrical activity produced by muscles, including, e.g., skeletal muscles. EMG may be performed using an instrument called an electromyograph, to produce a record called an electromyogram. An electromyograph may detect, among other things, the electrical potential generated by muscle cells when the cells are electrically or neurologically activated. The detected signals may be analyzed to detect, among other things, medical abnormalities, activation level, recruitment order or to analyze the biomechanics of human or animal movement. EMG is exceedingly intrusive in that it uses the insertion of needles through the skin and into the muscles and the use of these needles to measure electrical potential.\nElectrical impedance myography (EIM) is a novel technological approach to effectively address these limitations. Unlike standard electrophysiological approaches, EIM is less directly dependent upon inherent electrical potential of muscle or nerve tissue. EIM is based on electrical bioimpedance. It measures the effect of tissue structure and properties on flow of extremely small, non-intrusive amounts of electrical current. Unlike standard bioimpedance approaches, however, measurements can be performed over small areas of muscle and incorporate sophisticated analytic tools. In EIM, electrical current, such as, e.g., high-frequency alternating current, may be applied to localized areas of muscle via electrodes (e.g., surface electrodes) and the consequent surface voltage patterns may be analyzed. Although data can be obtained with off-the-shelf bioimpedance devices, these devices are far from ideal in terms of providing useful data reliably, as discussed in more detail below.\nFIG. 62A illustrates the concepts underlying EIM. Electrical current (sinusoid “a”) is applied via two or more outer surface electrodes generating a voltage difference measured by the two or more inner electrodes (sinusoid “b”). The voltage may be proportional to tissue resistance (R). Myocyte membrane lipid bilayers are capacitive in nature (e.g., they briefly store and then release some or all of the stored charge) and so exhibit reactance (X), making the voltage sine wave out of phase with applied current wave. Reactance and resistance values may be combined to obtain the summary phase angle (θ) via the relationship θ=arctan(X/R). FIG. 62B shows changes seen in diseased or less than ideal muscle tissue. Here, presence of connective tissue, fat and reduced muscle mass, among other things, may increase measured resistance; muscle fiber atrophy and loss also results in reduced reactance (e.g., timing of voltage sinusoid is now only slightly shifted relative to current). Thus, phase angle, as well as the resistance and reactance may be used to measure, e.g., disease progression. As disease advances, reactance and phase angle may decrease whereas resistance may increase.\nTwo additional aspects to EIM may include:                a) strong frequency dependence of EIM data. Thus, performing EIM measurements across a range of frequencies may help to characterize tissue. FIG. 63 shows multifrequency data “d” from a normal subject and from a patient with advanced ALS (Emory U. stem cell study), which is denoted by “e”. Note major alteration in impedance parameters across the frequency spectrum.        b) electrical anisotropy—directional dependence of current flow. Typically, electrical current flows relatively easily along muscle fibers than across them conferring a readily detectable anisotropy. Alteration in electrical anisotropy can also be used as a measure to evaluate muscle tissue to, e.g., determine a disease state, and early data show that anisotropy increases in, among other things, ALS.        \nAlthough much previous EIM work was done with off-the-shelf whole-body bioimpedance systems, for example, using these systems for localized impedance measurements may be problematic for a variety of reasons, including, but not limited to, the systems: 1) may not be calibrated for the very different impedances found in localized areas of tissue, such as, e.g., muscle tissue; 2) may be unable to effectively measure and account for muscle anisotropy; 3) rely on multiple, clumsy adhesive electrodes that may be slow to apply and result in spacing variability; and 4) may operate over a limited frequency range that may miss certain clinical information. Thus, there is a need for a handheld, rapidly applied, broadly capable, robust EIM system for bedside use.\nThere are some reports of the use of electrical impedance for biometric purposes. Examples of such uses may be found in: U.S. Pat. No. 6,122,544 to L. W. Orgon “Electrical Impedance Method and Apparatus for Detecting and Diagnosing Diseases” (Orgon 544); U.S. Pat. No. 6,768,921 to Leslie W. Organ, K. C. Smith, Reza Safaee-Rad, M. Graovac, G. P. Darmos, and I. Gavrilov, “Electrical impedance method and apparatus for detecting and diagnosing diseases” (Organ 921); U.S. Pat. No. 6,845,264 and PCT Application Publication No. WO 00/19894, Skladnev; Victor, Thompson; Richard L., Bath; Andrew R., “Apparatus for recognizing tissue types”, (Skladnev 264); U.S. Pat. No. 6,723,049 and Australian Application No. PR5718, Skladnev; Victor Nickolaevich, Blunsden; Christopher Kingsley, Stella; Rita “Apparatus for tissue type recognition using multiple measurement techniques” (Skladev 049); U.S. Pat. No. 7,212,852 to K. C. Smith, J. S. Ironstone, F. Zhang, “Bioimpedance measurement using controller-switched current injection and multiplexer selected electrode connection”, (Smith 852); U.S. Pat. No. 7,457,660 to K. C. Smith and J. I. Ironstone “Eliminating interface artifact errors in bioimpedance measurements” (Smith 660); U.S. Pat. No. 7,136,697 to Michael G. Singer “Methods for determining illness, progression to death, and/or timing of death of biological entity” (Singer 697); U.S. Pat. No. 7,003,346 to Michael G. Singer, “Method for illness and disease determination and management” (Singer 346); U.S. Pat. No. 8,103,337 to M. Gravovac, J I Marteus, Z. Pavlovic and J. Ironstone “Weighted Gradient Method and System for Diagnosing Disease” (Gravovac 337); U.S. Pat. No. 6,631,292 to R. J. Liedtke, (Liedtke 292) “Bio-electrical Impedance Analyzer”; U.S. Pat. No. 8,004,291 to Naosumi Waki, “Bioelectric impedance measuring circuit”, (Waki 291); U.S. Pat. No. 7,869,866 to Giannicola Loriga; Andrea Scozzari, “Device for the monitoring of physiologic variables through measurement of body electrical impedance”, (Loriga 866); U.S. Pat. No. 7,148,701 to Sin-Chong Park; In-Duk Hwang; “Apparatus for measuring electrical impedance” (Park); U.S. Patent Application Publication No. 2010/0292603 and PCT Application Publication No. WO/2007/035887 to C. A. Shiffman, R. Aaron and S. Rutkove, “Electrical Impedance Myography” (Shiffman 887); PCT Application No. WO 2011/022068 to Seward Rutkove “A Hand-held Device for Electrical Impedance Myography” (Rutkove 068); U.S. Pat. No. 5,919,142 and PCT Application No. PCT/GB96/01499 to Boone, Kevin Graham; Holder David Simon “Electrical impedance tomography method and apparatus (Boone 142); U.S. Patent Publication No. 2005/0004490 A1 to L. W. Organ, K. C. Smith, R Safaee-Rad, M. Granvac, P. Darmos and I Gavrilov, “Electrical Impedance Method and Apparatus for Detecting and Diagnosing Diseases” (Organ 490); U.S. Patent Application Publication No. 2005/0197591 to Z. Pavlovic, M Graovuc, J. S. Ironstone, “System and Method for Prebalancing Electrical Properties to Diagnose Disease” (Pavlovic 591); U.S. Patent Application Publication No. 2004/0073131 to L. W. Organ, K. C. Smith, R Sufaee-Rad, M. Graovac, G. P. Darmos and I. Gavrilov, “Electrical Impedance Method and Apparatus for Detecting and Diagnosing Diseases” (Organ 131); U.S. Patent Application Publication No. 2004/0167422 to L. W. Organ, R Sufaee-Rad, M. Graovac, K. C. Smith, J. S. Ironstone, “Breast Electrode Array and Method of Analysis for Detecting and Diagnosing Diseases”, (Organ 422); U.S. Patent Application Publication No. 2004/0210157 L. W. Organ, K. C. Smith, R. Safaee-Rad, M. Graovac, G. P. Darmos and I. Gavrilov “Electrical Impedance Method and Apparatus for Detecting and Diagnosing Diseases” (Organ 157); U.S. Patent Application Publication No. 2004/0210158 to L. W. Organ, K. C. Smith, R. Safaee-Rad, M. Graovac, G. P. Darmos and I. Gavrilov “Electrical Impedance Method and Apparatus for Detecting and Diagnosing Diseases” (Organ 158); U.S. Patent Application Publication No. 2004/0243018 to L. W. Organ, K. C. Smith and J. S. Ironstone, “Apparatus and Method for Determining Adequacy of Electrode-So-Skin Contact and Electrode Quality for Bioelectrical measurements” (Organ 018); U.S. Patent Application Publication No. 2004/0243019 to M. Graovac and Z. Pavlovic, “Weighted Gradient Method and System for Diagnosing Disease” (Graovac 019); U.S. Patent Application Publication No. 2008/0064979 to Z. Pavlovic, M. Graovac and J. S. Ironstone, “System and Method for Prebalancing Electrical Properties to Diagnose Disease” (Pavlovic 979); U.S. Patent Application Publication No. 2008/0076889 to L. W. Organ, R. Safaee-Rad, M. Graovac, K. C. Smith, J. S. Ironstone, “Breast Electrode Array and Method of Analysis for Detecting and Diagnosing Diseases”, (Organ 889); and U.S Patent Application Publication No. 2008/0249432 and PCT Application No. PCT/CA04/00458 A to Semlyen and M. Graovac “Diagnosis of Disease by Determination of Electrical Network Properties of a Body Part”, (Semiyen 432). All of these patents and patent applications are incorporated herein in their entirety by reference."}
{"text": "Ventilation systems commonly used in buildings, especially in spaces such as bedrooms and bathrooms, often include a ventilation duct to one end of which a fan is coupled. A ventilation device is arranged at the other end. Often, also one or more additional ventilation devices is/are connected to mouths of the ventilation ducts at different positions along the ventilation ducts. Said ventilation duct often extends over several different spaces in the property, for ventilation of those spaces. The ventilation device has an adjustable opening, here referred to as airflow opening, with which the airflow through the ventilation device between the ventilation duct and outside space can be adjusted. When the ventilation device is connected to a mouth of the ventilation duct, the airflow into or out of the ventilation duct can be adjusted by adjusting the size of the airflow opening.\nThe airflow through a ventilation device depends on factors such as the fan's effect, the dimensions of the ventilation duct, and the size of the airflow opening of the ventilation device. Here, the dimension of the ventilation duct refers to its diameter. When the ventilation system comprises a plurality of ventilation devices, such devices are generally set such that the various ventilation devices have different sizes of the airflow opening in order to thereby adjust the pressure distribution in the ventilation system. By adjusting the airflow opening of the various ventilation devices, unnecessarily high pressures can be throttled down. In this way, a predetermined airflow is obtainable through the respective ventilation devices, i.e. a desired degree of ventilation is obtainable in all spaces in which one or more ventilation devices is/are arranged. Too low airflow causes inadequate ventilation, while too high airflow causes increased energy costs.\nThe airflow, i.e. amount of feed air or exhaust air, is generally set according to current practice, in accordance with the dimensions of the ventilation duct.\nTo achieve that airflow calls for a certain pressure distribution in the ventilation system.\nIt is a problem with these systems that they generate sound which may be perceived as disturbing. For these ventilation systems there are therefore threshold values for maximum recommended sound effect level. Especially, sound is generated in the ventilation device at airflow through its opening towards the surroundings, i.e. the airflow opening. The threshold values for permissible sound effect level generated by the respective ventilation devices set limits to how large a pressure drop can be accomplished over the ventilation device, i.e. which opening degree the respective ventilation devices can have. This also sets limits to which airflow can be obtained through the ventilation device.\nAs mentioned above, a ventilation system usually comprises a plurality of ventilation devices at different distances from the fan. As the pressure generated by the fan is lowest at the ventilation device located furthest from the fan, said ventilation device is set with maximal opening, i.e. said ventilation device has maximal airflow opening size. The pressure required over said ventilation device to achieve a specified airflow determines the fan's operating conditions. To minimize energy consumption, a pressure drop should be as low as possible.\nAt the same time, a specified airflow must be obtained also through other ventilation devices being positioned closer to the fan and thus experiencing a higher pressure from the fan. Therefore, a certain degree of throttling of the pressure over the respective ventilation devices, a certain degree of pressure drop, is called for so that the specified airflow is neither exceeded nor underpassed. However, the recommended maximum sound effect level sets limits as to how much the pressure over a ventilation device can be throttled, because of the sounds occurring at airflow through the ventilation device. As will be described in more detail below, factors, such as the size of the airflow opening of the ventilation device, the dimensions of the ventilation device, and the size of an airflow there through, affect the sound effect level generated in the ventilation device at airflow there through. Therefore, the maximal degree of throttling of the pressure in a ventilation device obtainable over a ventilation device, without the recommended maximum sound effect level being exceeded, should be as high as possible to obtain effective ventilation throughout the entire ventilation system. Overall, these factors thus set limits as to the ventilation system.\nA ventilation system for feed air ventilation has been described above. The same applies also to exhaust air ventilation."}
{"text": "Hybrid circuits based on copper and compatible resistors should be encapsulated to insure resistor durability in humid atmospheres. Furthermore, manufacturers prefer glass encapsulation to protect the copper from long-term corrosion. Only contact pads are soldered.\nThe encapsulant system must exhibit several features which are difficult to achieve together. It must form a bubble-free seal at low enough firing temperature to prevent shift of the underlying resistors. If the glass flows too much, it will diffuse into the resistor and shift the value upward. If it does not flow enough, it will not seal. The organic vehicle necessary for screen printing must burn out in nitrogen at this low temperature and the surrounding copper must remain stain-free to allow solder acceptance. Copper will form an oxide stain at low temperature which will disappear when fired at higher temperature (.about.600.degree. C.) in nitrogen.\nThus an ideal encapsulant should screen print smoothly and rapidly with a vehicle which is decomposable in a low oxygen-containing atmosphere such as nitrogen at a low enough temperature to allow the glass to flow sufficiently to form a seal, but not so much as to shift the resistors. Furthermore, the firing temperature should be high enough to remove any oxide which forms on the copper conductors.\nA previous attempt to combine these features was the subject of U.S. patent application Ser. No. 12,133, filed Feb. 6, 1987 in which Pb(NO.sub.3).sub.2 was used to coat low melting glass and burn out the vehicle at low firing temperature. The low melting glass formed a good seal without excessive resistor shift. A problem with this approach is, however, that excess oxidizing agent [Pb(NO.sub.3).sub.2 ] required for complete vehicle removal and bubble-free film formation may also oxidize the copper and thus form an unsolderable stain. On the other hand, higher temperature firing which would remove the stain causes excess glass flow and consequently resistor shift."}
{"text": "Indoor positioning systems are useful for determining device location within buildings. One type of indoor positioning system is based on wireless access point fingerprinting. A common indoor wireless network is a WiFi network based on the IEEE 802.11 standard. An indoor positioning system stores “fingerprints” of expected wireless access point signal associated with various points or location in the facility.\nMeeting room allocation, scheduling and occupancy often relies upon a meeting organizer to correctly book the meeting room via a scheduling application and to faithfully adhere to start and end times. Many times people may use a meeting room without pre-booking the room via the scheduling system, or a pre-booked meeting may run late. Someone that has a subsequent booking of the meeting room may be unsure of whether the meeting room is occupied and whether he or she should possibly interrupt the meeting. If the meeting room is not equipped with a dedicated room occupancy sensor, e.g. a motion sensor, temperature sensor, or pressure sensor, for example, designed to detect the actual presence of a person, then someone outside the room may be unaware of whether the room is actually occupied.\nLike reference numerals are used in the drawings to denote like elements and features."}
{"text": "1. Field of the Invention\nThe present invention relates to a lead frame and a resin-molded-type semiconductor device using the same.\n2. Description of the Related Art\nA resin-molded-type semiconductor has been conventionally used mainly for low-frequency purposes. On the other hand, for high-frequency purposes a ceramic package has been generally and widely used. When a high frequency signal is transmitted, since it is difficult to obtain a desired characteristic impedance in a lead frame, the wave shapes of such high-frequency signals may easily be deformed. Therefore, for a high-frequency semiconductor device, a ceramic package can appropriately be used, because the ceramic package is able to be adapted to the desired characteristic impedance. In addition, a ceramic package having multi-layered line patterns can be formed and, therefore, the characteristic impedance thereof can appropriately be matched by suitably designing the distance between adjacent layers or line patterns or the thickness of respective layers.\nOn the other hand, a lead frame for resin-molded-type semiconductors has not been designed in such a manner that the width of leads or the distance between the adjacent leads are in conformity with those values to obtain a desired characteristic impedance. Therefore, a deflection of a signal may occur due to the miss-matching of the characteristic impedance and, therefore, decay of signal or deformation of the signal wave shape may easily occur. Thus, in general, a conventional lead frame is inappropriate for a high-frequency semiconductor device."}
{"text": "The systemic administration of therapeutic agents treats the body as a whole even though the disease to be treated may be localized. In some cases of localized disease, systemic administration may not be desirable because the drug agents may have unwanted effects on parts of the body which are not to be treated or because treatment of the diseased part of the body requires a high concentration of a drug agent that may not be achievable by systemic administration.\nIt is therefore often desirable to administer therapeutic agents to only localized sites within the body. Common examples of where this is needed include cases of localized disease (e.g., coronary heart disease) and occlusions, lesions, or other disease in body lumens. Several devices and methods for localized drug delivery are known. In one example, such devices are drug delivery balloons, and methods of their use include the steps of coating a balloon attached to a balloon catheter with a drug and a carrier matrix, inserting the catheter into a blood vessel, tracking the balloon to a desired location, and expanding the balloon against the surrounding tissue to transfer the drug locally at the intended treatment site.\nOne of the potential drawbacks to localized drug delivery is the possibility of premature or unintended release of the drug, the carrier matrix, and/or the drug/carrier matrix combination. This may occur during tracking and placement at the treatment site of a drug delivery device and post delivery as the device is withdrawn from the body. Such unintended release may result from drug diffusion, device contact with areas proximate the treatment site, or washing of the drug from the surface of the delivery device due to blood flow. This is of particular concern when the device comprises a therapeutic agent of a type or dosage not intended to be released to tissue or blood outside the treatment site.\nDrugs or coating components shed in this unwanted fashion may be in particulate form or may be in solution. The downstream release of undesirable particles is known as “particulation”. For example, particulation of large particles can create problems such as ischemia in tissues, especially in tissues supplied by small diameter vessels. Furthermore, the resulting effects of biodistribution of such particles are not well understood and may result in adverse effects.\nIn view of the potential drawbacks to current, localized drug delivery, there exists a need for devices and methods that allow for controlled, localized delivery of drug agents to specific treatment sites within a mammalian body that avoids premature or unintended drug release away from the intended treatment site."}
{"text": "The present invention relates generally to wearable computers and, more specifically, to methods and arrangements for facilitating the connection of auxiliary components therewith.\nAs miniaturization of computer products continues, the problem of adequate space for connections of input and output devices (I/O) continues to grow. The minimum size of many computer products are now being defined by the number and type of connectors used, especially on the highly integrated devices of the type used in wearable computer technology. In other words, since it is the case that a minimum amount of space must typically be allocated for accommodating connectors or connecting interfaces in the environment of a wearable computer, it has been found that further efforts towards miniaturization in the context of wearable computers has been hampered.\nA need has thus been recognized in connection with facilitating further miniaturization in the context of wearable computers by reducing the amount of space associated with connectors, connection interfaces or other connecting elements.\nIn accordance with at least one presently preferred embodiment of the present invention, a single docking-type connector is provided on a computer system unit, and distributes the I/O connectors required for a given task, into the harness (wearing device) of the computer. The docking connector preferably has enough pins to carry all of the I/O signals, power and grounds for the computer.\nSuccessful wearable computers will offer flexible configurations of I/O and wearing possibilities. Because the system unit has all of the possible I/O signals, different types of harnesses (as variously illustrated by way of example in the drawings) allow for many configurations. The harnesses with the cabling are relatively inexpensive compared to the system unit.\nGenerally, in accordance with at least one presently preferred embodiment of the present invention, problems relating to the space taken up by connectors are solved by essentially exchanging the multiple connectors required for I/O on a system unit for a single docking connector. By redistributing the individual connectors into various harnesses and docking stations, the system unit is greatly reduced in size and can achieve. attendant advantages as discussed herein.\nIn one aspect, the present invention provides a connector system for a computer, the connector system comprising: a docking connector which interfaces with a system unit; a plurality of leads in communication with the docking connector; and a plurality of connecting elements associated with the plurality of leads, the connecting elements each being adapted to interface with at least one external component.\nIn another aspect, the present invention provides a harness system for supporting a portable computer, the harness system comprising: a harness which accommodates a system unit; a docking connector which interfaces with the system unit; a plurality of leads in communication with the docking connector; and a plurality of connecting elements associated with the plurality of leads, the connecting elements each being adapted to interface with at least one external component.\nFor a better understanding of the present invention, together with other and further features and advantages thereof, reference is made to the following description, taken in conjunction with the accompanying drawings, and the scope of the invention will be pointed out in the appended claims."}
{"text": "1. Field of the Invention\nThis invention relates generally to CRF receptor antagonists, and to methods of treating disorders by administration of such antagonists to a warm-blooded animal in need thereof.\n2. Description of the Related Art\nThe first corticotropin-releasing factor (CRF) was isolated from ovine hypothalmi and identified as a 41-amino acid peptide (Vale et al., Science 213:1394-1397, 1981). Subsequently, sequences of human and rat CRF were isolated and determined to be identical, but different from ovine CRF in 7 of the 41 amino acid residues (Rivier et al., Proc. Natl. Acad. Sci. USA 80:4851, 1983; Shibahara et al., EMBO J. 2:775, 1983).\nCRF has been found to produce profound alterations in endocrine, nervous and immune system function. CRF is believed to be the major physiological regulator of the basal and stress-release of adrenocorticotropic hormone (xe2x80x9cACTHxe2x80x9d), xcex1-endorphin, and other pro-opiomelanocortin (xe2x80x9cPOMCxe2x80x9d)-derived peptides from the anterior pituitary (Vale et al., Science 213:1394-1397, 1981). Briefly, CRF is believed to initiate its biological effects by binding to a plasma membrane receptor which has been found to be distributed throughout the brain (DeSouza et al., Science 224:1449-1451, 1984), pituitary (DeSouza et al., Methods Enzymol. 124:560, 1986; Wynn et al., Biochem. Biophys. Res. Comm. 110:602-608, 1983), adrenals (Udelsman et al., Nature 319:147-150, 1986) and spleen (Webster, E. L., and E. B. DeSouza, Endocrinology 122:609-617, 1988). The CRF receptor is coupled to a GTP-binding protein (Perrin et al., Endocrinology 118:1171-1179, 1986) which mediates CRF-stimulated increase in intracellular production of cAMP (Bilezikjian, L. M., and W. W. Vale, Endocrinology 113:657-662, 1983). The receptor for CRF has now been cloned from rat (Perrin et al., Endo 133(6):3058-3061, 1993), and human brain (Chen et al., PNAS 90(19):8967-8971, 1993; Vita et al., FEBS 335(1):1-5, 1993). This receptor is a 415 amino acid protein comprising seven membrane spanning domains. A comparison of identity between rat and human sequences shows a high degree of homology (97%) at the amino acid level.\nIn addition to its role in stimulating the production of ACTH and POMC, CRF is also believed to coordinate many of the endocrine, autonomic, and behavioral responses to stress, and may be involved in the pathophysiology of affective disorders. Moreover, CRF is believed to be a key intermediary in communication between the immune, central nervous, endocrine and cardiovascular systems (Crofford et al., J. Clin. Invest. 90:2555-2564, 1992; Sapolsky et al., Science 238:522-524, 1987; Tilders et al., Regul. Peptides 5:77-84, 1982). Overall, CRF appears to be one of the pivotal central nervous system neurotransmitters and plays a crucial role in integrating the body\"\"s overall response to stress.\nAdministration of CRF directly to the brain elicits behavioral, physiological, and endocrine responses identical to those observed for an animal exposed to a stressful environment. For example, intracerebroventricular injection of CRF results in behavioral activation (Sutton et al., Nature 297:331, 1982), persistent activation of the electroencephalogram (Ehlers et al., Brain Res. 278:332, 1983), stimulation of the sympathoadrenomedullary pathway (Brown et al., Endocrinology 110:928, 1982), an increase of heart rate and blood pressure (Fisher et al., Endocrinology 110:2222, 1982), an increase in oxygen consumption (Brown et al., Life Sciences 30:207, 1982), alteration of gastrointestinal activity (Williams et al., Am. J. Physiol. 253:G582, 1987), suppression of food consumption (Levine et al., Neuropharmacology 22:337, 1983), modification of sexual behavior (Sirinathsinghji et al., Nature 305:232, 1983), and immune function compromise (Irwin et al., Am. J. Physiol. 255:R744, 1988). Furthermore, clinical data suggests that CRF may be hypersecreted in the brain in depression, anxiety-related disorders, and anorexia nervosa. (DeSouza, Ann. Reports in Med. Chem. 25:215-223, 1990). Accordingly, clinical data suggests that CRF receptor antagonists may represent novel antidepressant and/or anxiolytic drugs that may be useful in the treatment of the neuropsychiatric disorders manifesting hypersecretion of CRF.\nThe first CRF receptor antagonists were peptides (see, e.g., Rivier et al., U.S. Pat. No. 4,605,642; Rivier et al., Science 224:889, 1984). While these peptides established that CRF receptor antagonists can attenuate the pharmacological responses to CRF, peptide CRF receptor antagonists suffer from the usual drawbacks of peptide therapeutics including lack of stability and limited oral activity. More recently, small molecule CRF receptor antagonists have been reported. For example, substituted 4-thio-5-oxo-3-pyyrazoline derivatives (Abreu et al., U.S. Pat. No. 5,063,245) and substituted 2-aminothiazole derivatives (Courtemanche et al., Australian Patent No. AU-A-41399/93) have been reported as CRF receptor antagonists. These particular derivatives were found to be effective in inhibiting the binding of CRF to its receptor in the 1-10 xcexcM range and 0.1-10 xcexcM range, respectively.\nDue to the physiological significance of CRF, the development of biologically-active small molecules having significant CRF receptor binding activity and which are capable of antagonizing the CRF receptor remains a desirable goal. Such CRF receptor antagonists would be useful in the treatment of endocrine, psychiatric and neurologic conditions or illnesses, including stress-related disorders in general.\nWhile significant strides have been made toward achieving CRF regulation through administration of CRF receptor antagonists, there remains a need in the art for effective small molecule CRF receptor antagonists. There is also a need for pharmaceutical compositions containing such CRF receptor antagonists, as well as methods relating to the use thereof to treat, for example, stress-related disorders. The present invention fulfills these needs, and provides other related advantages.\nIn brief, this invention is generally directed to CRF receptor antagonists, and more specifically to CRF receptor antagonists having the following general structure (I): \nincluding stereoisomers and pharmaceutically acceptable salts thereof, wherein X, R1, R2 and Ar are as defined below.\nThe CRF receptor antagonists of this invention have utility over a wide range of therapeutic applications, and may be used to treat a variety of disorders or illnesses, including stress-related disorders. Such methods include administering an effective amount of a CRF receptor antagonist of this invention, preferably in the form of a pharmaceutical composition, to an animal in need thereof. Accordingly, in another embodiment, pharmaceutical compositions are disclosed containing one or more CRF receptor antagonists of this invention in combination with a pharmaceutically acceptable carrier and/or diluent.\nThese and other aspects of the invention will be apparent upon reference to the following detailed description. To this end, various references are set forth herein which describe in more detail certain procedures, compounds and/or compositions, and are hereby incorporated by reference in their entirety.\nThe present invention is directed generally to compounds useful as corticotropin-releasing factor (CRF) receptor antagonists.\nIn a first embodiment, the CRF receptor antagonists of this invention have the following structure (I): \nincluding stereoisomers and pharmaceutically acceptable salts thereof,\nwherein: \nrepresents xe2x80x94Nxe2x95x90CHxe2x80x94, xe2x80x94NHxe2x80x94CH2xe2x80x94 or xe2x80x94NHxe2x80x94(CH2)2xe2x80x94;\nX is N or CR3;\nR1 is xe2x80x94CH(R4)(R5);\nR2 is C1-6alkyl;\nR3 is hydrogen or C1-6alkyl;\nR4 is hydrogen, C1-6alkyl, mono- or di(C3-6cycloalkyl)methyl, C3-6cycloalkyl, C3-6alkenyl, hydroxyC1-6alkyl, C1-6alkylcarbonyloxyC1-6alkyl, or C1-6alkyloxyC1-6alkly, and\nR5 is C1-8alkyl, mono- or di(C3-6cycloalkyl)methyl, Ar1CH2, C3-6alkenyl, C1-6alkyloxyC1-6alkyl, hydroxyC1-6alkyl, thienylmethyl, furanylmethyl, C1-6alkylthioC1-6alkyl, morpholinyl, mono- or di(C1-6alkyl)aminoC1-6alkyl, di(C1-6alkyl)amino, C1-6alkylcarbonylC1-6alkyl, C1-6alkyl substituted with imidazolyl, or a radical of the formula xe2x80x94(C1-6alkanediyl)xe2x80x94Oxe2x80x94COxe2x80x94Ar1,\nor R4 and R5 taken together with the carbon atom to which they are bonded form a C5-8cycloalkyl optionally substituted with one or more substituents independently selected from C1-6alkyl;\nAr is phenyl substituted with 1, 2 or 3 substituents independently selected from halo, C1-6alkyl, trifluoromethyl, cyano, C1-6alkyloxy, benzyloxy, C1-6alkylthio, nitro, amino, and mono- or di(C1-6alkyl)amino; or an aromatic C3-12heterocycle optionally substituted with 1, 2 or 3 substituents independently selected from halo, C1-6alkyl, trifluoromethyl, hydroxy, cyano, C1-6alkyloxy, benzyloxy, C1-6alkylthio, nitro, amino, mono- or di(C1-6alkyl)amino, and piperidinyl; and\nAr1 is phenyl, pyridinyl, or phenyl substituted with 1, 2 or 3 substituents independently selected from halo, C1-6alkyl, C1-6alkyloxy, di(C1-6alkyl)aminoC1-6alkyl, trifluoromethyl and C1-6alkyl substituted with morpholinyl.\nIn the context of this invention, the preceding terms have the meanings set forth below.\nxe2x80x9cC1-6alkylxe2x80x9d or xe2x80x9cC1-8alkylxe2x80x9d represents a straight chain or branched alkyl having from 1 to 6 carbon atoms or 1 to 8 carbon atoms, respectively, such as methyl, ethyl, n-propyl, isopropyl, n-butyl, tert-butyl, n-pentyl, and the like.\nxe2x80x9cC1-6alkyloxyxe2x80x9d represents the group xe2x80x94O(C1-6alkyl).\nxe2x80x9cC1-6alkylthioxe2x80x9d represents the group xe2x80x94S(C1-6alkyl).\nxe2x80x9cC3-6cycloalkylxe2x80x9d represents a cyclic alkyl having from 3 to 6 carbon atoms, including cyclopropyl, cyclopentyl, cyclopentyl, and cyclohexyl.\nxe2x80x9cC5-8cycloalkylxe2x80x9d represents a cyclic alkyl having from 5 to 8 carbon atoms, such as cyclopentyl, cyclohexyl, and the like.\nxe2x80x9cC3-6alkenylxe2x80x9d represents an unsaturated straight chain or branched alkyl having from 3 to 6 carbon atoms, and having at least one double bond, such as propylene, 1-butene, 2-butene, 2-methylpropene, and the like.\nxe2x80x9cHydroxyC1-6alkylxe2x80x9d represents a C1-6alkyl substituted with at least one hydroxyl group.\nxe2x80x9cMono- or di(C3-6cycloalkyl)methylxe2x80x9d represents a methyl group substituted with one or two C3-6cycloalkyl groups, such as cyclopropylmethyl, dicyclopropylmethyl, and the like.\nxe2x80x9cC1-6alkylcarbonylC1-6alkylxe2x80x9d represents a C1-6alkyl substituted with a xe2x80x94COC1-6alkyl group.\nxe2x80x9cC1-6alkylcarbonyloxyC1-6alkylxe2x80x9d represents a C1-6alkyl substituted with a xe2x80x94COOC1-6alkyl group.\nxe2x80x9cC1-6alkyloxyC1-6alkylxe2x80x9d represents a C1-6alkyl substituted with a xe2x80x94OC1-6alkyl group.\nxe2x80x9cC1-6alkylthioC1-6alkylxe2x80x9d represents a C1-6alkyl substituted with a xe2x80x94SC1-6alkyl group.\nxe2x80x9cMono- or di(C1-6alkyl)amino represents an amino substituted with one C1-6alkyl or with two C1-6alkyls, respectively.\nxe2x80x9cMono- or di(C1-6alkyl)aminoC1-6alkylxe2x80x9d represents a C1-6alkyl substituted with a mono- or di(C1-6alkyl)amino.\nxe2x80x9cC1-6alkanediylxe2x80x9d represents a divalent C1-6alkyl radical, such as methylene (xe2x80x94CH2xe2x80x94), ethylene (xe2x80x94CH2CH2xe2x80x94), and the like.\nxe2x80x9cC3-12heterocyclexe2x80x9d represents a ring made up of more than one kind of atom, and which contains 3 to 12 carbon atoms, such as pyridinyl, pyrimidinyl, furanyl, thienyl, imidazolyl, thiazolyl, pyrazolyl, pyridazinyl, pyrazinyl, triazinyl (such as 1,3,5), and the like.\nxe2x80x9cHaloxe2x80x9d means fluoro, chloro, bromo or iodo.\nRepresentative CRF receptor antagonists of this invention include compounds having the following structures (Ia), (Ib) and (Ic): \nWhen X of compounds (Ia), (Ib) and (Ic) is N, representative compounds of this invention include the following compounds (Iaxe2x80x2), (Ibxe2x80x2) and (Icxe2x80x2), and when X is CR3, and R3 is hydrogen, representative compounds of this invention include the following compounds (Iaxe2x80x3), (Ibxe2x80x3) and (Icxe2x80x3): \nIn the embodiment where the R4 and R5 groups of R1 taken together form a C5-8cycloalkyl, the resulting R1 group has the structure: \nWhen the above structure is optionally substituted with one or more C1-6alkyl groups, a representative R1 moiety has the following structure: \nwherein R6 and R7 are the same or different and independently selected from a C1-6alkyl, such as methyl or ethyl.\nRepresentative Ar, R1 and R2 groups of this invention are set forth in the following Table. To this end, it should be understood that each combination of the Ar, R1 and R2 groups listed in the following Table 1 represents individual compounds of structure (I), as well as the more specific structures (Ia), (Ib) and (Ic) and sub-structures (Iaxe2x80x2), (Iaxe2x80x3), (Ibxe2x80x2), (Ibxe2x80x3), (Icxe2x80x2) and (Icxe2x80x3).\nThe compounds of the present invention may be prepared by known organic synthesis techniques, including the methods described in more detail in the Examples, and may generally be utilized as the free base. Alternatively, the compounds of this invention may be used in the form of acid addition salts. Acid addition salts of the free base amino compounds of the present invention may be prepared by methods well known in the art, and may be formed from organic and inorganic acids. Suitable organic acids include maleic, fumaric, benzoic, ascorbic, succinic, methanesulfonic, acetic, oxalic, propionic, tartaric, salicylic, citric, gluconic, lactic, mandelic, cinnamic, aspartic, stearic, palmitic, glycolic, glutamic, and benzenesulfonic acids. Suitable inorganic acids include hydrochloric, hydrobromic, sulfuric, phosphoric, and nitric acids.\nThe effectiveness of a compound as a CRF receptor antagonist may be determined by various assay methods. Suitable CRF antagonists of this invention are capable of inhibiting the specific binding of CRF to its receptor and antagonizing activities associated with CRF. A compound of structure (I) may be assessed for activity as a CRF antagonist by one or more generally accepted assays for this purpose, including (but not limited to) the assays disclosed by DeSouza et al. (J. Neuroscience 7:88, 1987) and Battaglia et al. (Synapse 1:572, 1987). As mentioned above, suitable CRF antagonists include compounds which demonstrate CRF receptor affinity. CRF receptor affinity may be determined by binding studies that measure the ability of a compound to inhibit the binding of a radiolabeled CRF (e.g., [125I]tyrosine-CFR) to its receptor (e.g., receptors prepared from rat cerebral cortex membranes). The radioligand binding assay described by DeSouza et al. (supra, 1987) provides an assay for determining a compound\"\"s affinity for the CRF receptor. Such activity is typically calculated from the IC50 as the concentration of a compound necessary to displace 50% of the radiolabeled ligand from the receptor, and is reported as a xe2x80x9cKixe2x80x9d value calculated by the following equation:       K    i    =            IC      50              1      +              L        /                  K          D                    \nwhere L=radioligand and KD=affinity of radioligand for receptor (Cheng and Prusoff, Biochem. Pharmacol. 22:3099, 1973).\nIn addition to inhibiting CRF receptor binding, a compound\"\"s CRF receptor antagonist activity may be established by the ability of the compound to antagonize an activity associated with CRF. For example, CRF is known to stimulate various biochemical processes, including adenylate cyclase activity. Therefore, compounds may be evaluated as CRF antagonists by their ability to antagonize CRF-stimulated adenylate cyclase activity by, for example, measuring cAMP levels. The CRF-stimulated adenylate cyclase activity assay described by Battaglia et al. (supra, 1987) provides an assay for determining a compound\"\"s ability to antagonize CRF activity. Accordingly, CRF receptor antagonist activity may be determined by assay techniques which generally include an initial binding assay (such as disclosed by DeSouza (supra, 1987)) followed by a cAMP screening protocol (such as disclosed by Battaglia (supra, 1987)).\nWith reference to CRF receptor binding affinities, CRF receptor antagonists of this invention have a Ki of less than 10 xcexcM. In a preferred embodiment of this invention, a CRF receptor antagonist has a Ki of less than 1 xcexcM, and more preferably less than 0.25 xcexcM (i.e., 250 nM).\nThe CRF receptor antagonists of the present invention demonstrate activity at the CRF receptor site, and may be used as therapeutic agents for the treatment of a wide range of disorders or illnesses including endocrine, psychiatric, and neurologic disorders or illnesses. More specifically, the CRF receptor antagonists of the present invention may be useful in treating physiological conditions or disorders arising from the hypersecretion of CRF. Because CRF is believed to be a pivotal neurotransmitter that activates and coordinates the endocrine, behavioral and automatic responses to stress, the CRF receptor antagonists of the present invention can be used to treat neuropsychiatric disorders. Neuropsychiatric disorders which may be treatable by the CRF receptor antagonists of this invention include affective disorders such as depression; anxiety-related disorders such as generalized anxiety disorder, panic disorder, obsessive-compulsive disorder, abnormal aggression, cardiovascular abnormalities such as unstable angina and reactive hypertension; and feeding disorders such as anorexia nervosa, bulimia, and irritable bowel syndrome. CRF antagonists may also be useful in treating stress-induced immune suppression associated with various disease states, as well as stroke. Other uses of the CRF antagonists of this invention include treatment of inflammatory conditions (such as rheumatoid arthritis, uveitis, asthma, inflammatory bowel disease and G.I. motility), Cushing\"\"s disease, infantile spasms, epilepsy and other seizures in both infants and adults, and various substance abuse and withdrawal (including alcoholism).\nIn another embodiment, compounds of this invention and their analogs may be used as Positron Emission Tomography (PET) ligands, Single Photon Emission Computed Tomography (SPECT) ligands, or other diagnostic radiopharmaceutical agents. Incorporation of an appropriate isotope (such as 11C or 18F for PET or 125I in the case of SPECT) may provide an agent useful for the diagnosis or therapeutic management of a patient. In addition, use of a compound of the present invention may provide a physiological, functional, or biological assessment of a patient or provide disease or pathology detection and assessment.\nIn another embodiment of the invention, pharmaceutical compositions containing one or more CRF receptor antagonists are disclosed. For the purposes of administration, the compounds of the present invention may be formulated as pharmaceutical compositions. Pharmaceutical compositions of the present invention comprise a CRF receptor antagonist of the present invention (i.e., a compound of structure (I)) and a pharmaceutically acceptable carrier and/or diluent. The CRF receptor antagonist is present in the composition in an amount which is effective to treat a particular disorderxe2x80x94that is, in an amount sufficient to achieve CRF receptor antagonist activity, and preferably with acceptable toxicity to the patient. Preferably, the pharmaceutical compositions of the present invention may include a CRF receptor antagonist in an amount from 0.1 mg to 250 mg per dosage depending upon the route of administration, and more preferably from 1 mg to 60 mg. Appropriate concentrations and dosages can be readily determined by one skilled in the art.\nPharmaceutically acceptable carrier and/or diluents are familiar to those skilled in the art. For compositions formulated as liquid solutions, acceptable carriers and/or diluents include saline and sterile water, and may optionally include antioxidants, buffers, bacteriostats and other common additives. The compositions can also be formulated as pills, capsules, granules, or tablets which contain, in addition to a CRF receptor antagonist, diluents, dispersing and surface active agents, binders, and lubricants. One skilled in this art may further formulate the CRF receptor antagonist in an appropriate manner, and in accordance with accepted practices, such as those disclosed in Remington\"\"s Pharmaceutical Sciences, Gennaro, Ed., Mack Publishing Co., Easton, Pa. 1990.\nIn another embodiment, the present invention provides a method for treating a variety of disorders or illnesses, including endocrine, psychiatric and neurologic disorders or illnesses. Such methods include administering of a compound of the present invention to a warm-blooded animal in an amount sufficient to treat the disorder or illness. Such methods include systemic administration of a CRF receptor antagonist of this invention, preferably in the form of a pharmaceutical composition. As used herein, systemic administration includes oral and parenteral methods of administration. For oral administration, suitable pharmaceutical compositions of CRF receptor antagonists include powders, granules, pills, tablets, and capsules as well as liquids, syrups, suspensions, and emulsions. These compositions may also include flavorants, preservatives, suspending, thickening and emulsifying agents, and other pharmaceutically acceptable additives. For parental administration, the compounds of the present invention can be prepared in aqueous injection solutions which may contain, in addition to the CRF receptor antagonist, buffers, antioxidants, bacteriostats, and other additives commonly employed in such solutions.\nAs mentioned above, administration of a compound of the present invention can be used to treat a wide variety of disorders or illnesses. In particular, the compounds of the present invention may be administered to a warm-blooded animal for the treatment of depression, anxiety disorder, panic disorder, obsessive-compulsive disorder, abnormal aggression, unstable angina, reactive hypertension, anorexia nervosa, bulimia, irritable bowel syndrome, stress-induced immune suppression, stroke, inflammation, Cushing\"\"s disease, infantile spasms, epilepsy, and substance abuse or withdrawal.\nThe following examples are provided for purposes of illustration, not limitation."}
{"text": "1. Field of the Invention\nThis invention relates to a roofing cover system wherein the exposed tab portions of old asphalt shingles and the like are covered by metal cover panels in a manner interlocking the cover panels to the roofing tabs and securing, with shank type fasteners, the cover panels and free lower corner portions of the old roofing tabs to the underlying roof structure.\n2. Description of Related Art\nVarious different forms of roofing cladding systems heretofore have been provided such as those disclosed in U.S. Pat. Nos. 632,691, 1,039,065, 1,090,531, 1,318,125, 2,038,192, 2,231,008 and 2,450,562. However, these previously known roofing covering systems do not include shingle tab covering panels which interlocking engage with each other as well as the underlying roofing tabs and which may be secured in operative position through the utilization of a pair of shank type fasteners secured downwardly through not the covering panels and the shingles of the tabs being covered but also the free lower corner portions of shingle tabs in the next above course of shingles."}
{"text": "This invention relates to particles for use in electrophoretic displays, and to electrophoretic media and displays incorporating such particles. More specifically, this invention relates to particles which enable the use, in electrophoretic media and displays, of organic pigments which, in their normal form, have particle sizes too small for use in electrophoretic media and displays.\nThe terms “bistable” and “bistability” are used herein in their conventional meaning in the art to refer to displays comprising display elements having first and second display states differing in at least one optical property, and such that after any given element has been driven, by means of an addressing pulse of finite duration, to assume either its first or second display state, after the addressing pulse has terminated, that state will persist for at least several times, for example at least four times, the minimum duration of the addressing pulse required to change the state of the display element. It is shown in U.S. Pat. No. 7,170,670 that some particle-based electrophoretic displays capable of gray scale are stable not only in their extreme black and white states but also in their intermediate gray states, and the same is true of some other types of electro-optic displays. This type of display is properly called “multi-stable” rather than bistable, although for convenience the term “bistable” may be used herein to cover both bistable and multi-stable displays.\nParticle-based electrophoretic displays have been the subject of intense research and development for a number of years. In this type of display, a plurality of charged particles move through a fluid under the influence of an electric field. Electrophoretic displays can have attributes of good brightness and contrast, wide viewing angles, state bistability, and low power consumption when compared with liquid crystal displays. In such electrophoretic displays, an optical property is changed by application of the electric field; this optical property is typically color perceptible to the human eye, but may be another optical property, such as optical transmission, reflectance, luminescence or, in the case of displays intended for machine reading, pseudo-color in the sense of a change in reflectance of electromagnetic wavelengths outside the visible range. Nevertheless, problems with the long-term image quality of these displays have prevented their widespread usage. For example, particles that make up electrophoretic displays tend to settle, resulting in inadequate service-life for these displays.\nAs noted above, electrophoretic media require the presence of a fluid. In most prior art electrophoretic media, this fluid is a liquid, but electrophoretic media can be produced using gaseous fluids; see, for example, Kitamura, T., et al., “Electrical toner movement for electronic paper-like display”, IDW Japan, 2001, Paper HCS1-1, and Yamaguchi, Y., et al., “Toner display using insulative particles charged triboelectrically”, IDW Japan, 2001, Paper AMD4-4). See also U.S. Patent Publication No. 2005/0001810; European Patent Applications 1,462,847; 1,482,354; 1,484,635; 1,500,971; 1,501,194; 1,536,271; 1,542,067; 1,577,702; 1,577,703; and 1,598,694; and International Applications WO 2004/090626; WO 2004/079442; and WO 2004/001498. Such gas-based electrophoretic media appear to be susceptible to the same types of problems due to particle settling as liquid-based electrophoretic media, when the media are used in an orientation which permits such settling, for example in a sign where the medium is disposed in a vertical plane. Indeed, particle settling appears to be a more serious problem in gas-based electrophoretic media than in liquid-based ones, since the lower viscosity of gaseous suspending fluids as compared with liquid ones allows more rapid settling of the electrophoretic particles.\nNumerous patents and applications assigned to or in the names of the Massachusetts Institute of Technology (MIT) and E Ink Corporation have recently been published describing encapsulated electrophoretic media. Such encapsulated media comprise numerous small capsules, each of which itself comprises an internal phase containing electrophoretically-mobile particles suspended in a liquid suspending medium, and a capsule wall surrounding the internal phase. Typically, the capsules are themselves held within a polymeric binder to form a coherent layer positioned between two electrodes. Encapsulated media of this type are described, for example, in U.S. Pat. Nos. 5,930,026; 5,961,804; 6,017,584; 6,067,185; 6,118,426; 6,120,588; 6,120,839; 6,124,851; 6,130,773; 6,130,774; 6,172,798; 6,177,921; 6,232,950; 6,249,271; 6,252,564; 6,262,706; 6,262,833; 6,300,932; 6,312,304; 6,312,971; 6,323,989; 6,327,072; 6,376,828; 6,377,387; 6,392,785; 6,392,786; 6,413,790; 6,422,687; 6,445,374; 6,445,489; 6,459,418; 6,473,072; 6,480,182; 6,498,114; 6,504,524; 6,506,438; 6,512,354; 6,515,649; 6,518,949; 6,521,489; 6,531,997; 6,535,197; 6,538,801; 6,545,291; 6,580,545; 6,639,578; 6,652,075; 6,657,772; 6,664,944; 6,680,725; 6,683,333; 6,704,133; 6,710,540; 6,721,083; 6,724,519; 6,727,881; 6,738,050; 6,750,473; 6,753,999; 6,816,147; 6,819,471; 6,822,782; 6,825,068; 6,825,829; 6,825,970; 6,831,769; 6,839,158; 6,842,167; 6,842,279; 6,842,657; 6,864,875; 6,865,010; 6,866,760; 6,870,661; 6,900,851; 6,922,276; 6,950,200; 6,958,848; 6,967,640; 6,982,178; 6,987,603; 6,995,550; 7,002,728; 7,012,600; 7,012,735; 7,023,420; 7,030,412; 7,030,854; 7,034,783; 7,038,655; 7,061,663; 7,071,913; 7,075,502; 7,075,703; 7,079,305; 7,106,296; 7,109,968; 7,110,163; 7,110,164; 7,116,318; 7,116,466; 7,119,759; 7,119,772; 7,148,128; 7,167,155; 7,170,670; 7,173,752; 7,176,880; 7,180,649; 7,190,008; 7,193,625; 7,202,847; 7,202,991; 7,206,119; 7,223,672; 7,230,750; 7,230,751; 7,236,790; and 7,236,792; and U.S. Patent Applications Publication Nos. 2002/0060321; 2002/0090980; 2003/0011560; 2003/0102858; 2003/0151702; 2003/0222315; 2004/0094422; 2004/0105036; 2004/0112750; 2004/0119681; 2004/0136048; 2004/0155857; 2004/0180476; 2004/0190114; 2004/0196215; 2004/0226820; 2004/0257635; 2004/0263947; 2005/0000813; 2005/0007336; 2005/0012980; 2005/0017944; 2005/0018273; 2005/0024353; 2005/0062714; 2005/0067656; 2005/0099672; 2005/0122284; 2005/0122306; 2005/0122563; 2005/0134554; 2005/0151709; 2005/0152018; 2005/0156340; 2005/0179642; 2005/0190137; 2005/0212747; 2005/0213191; 2005/0219184; 2005/0253777; 2005/0280626; 2006/0007527; 2006/0024437; 2006/0038772; 2006/0139308; 2006/0139310; 2006/0139311; 2006/0176267; 2006/0181492; 2006/0181504; 2006/0194619; 2006/0197736; 2006/0197737; 2006/0197738; 2006/0202949; 2006/0223282; 2006/0232531; 2006/0245038; 2006/0256425; 2006/0262060; 2006/0279527; 2006/0291034; 2007/0035532; 2007/0035808; 2007/0052757; 2007/0057908; 2007/0069247; 2007/0085818; 2007/0091417; 2007/0091418; 2007/0097489; 2007/0109219; 2007/0128352; and 2007/0146310; and International Applications Publication Nos. WO 00/38000; WO 00/36560; WO 00/67110; and WO 01/07961; and European Patents Nos. 1,099,207 B1; and 1,145,072 B1.\nMany of the aforementioned patents and applications recognize that the walls surrounding the discrete microcapsules in an encapsulated electrophoretic medium could be replaced by a continuous phase, thus producing a so-called polymer-dispersed electrophoretic display, in which the electrophoretic medium comprises a plurality of discrete droplets of an electrophoretic fluid and a continuous phase of a polymeric material, and that the discrete droplets of electrophoretic fluid within such a polymer-dispersed electrophoretic display may be regarded as capsules or microcapsules even though no discrete capsule membrane is associated with each individual droplet; see for example, the aforementioned U.S. Pat. No. 6,866,760. Accordingly, for purposes of the present application, such polymer-dispersed electrophoretic media are regarded as sub-species of encapsulated electrophoretic media.\nA related type of electrophoretic display is a so-called “microcell electrophoretic display”. In a microcell electrophoretic display, the charged particles and the fluid are not encapsulated within microcapsules but instead are retained within a plurality of cavities formed within a carrier medium, typically a polymeric film. See, for example, U.S. Pat. Nos. 6,672,921 and 6,788,449, both assigned to Sipix Imaging, Inc.\nAlthough electrophoretic media are often opaque (since, for example, in many electrophoretic media, the particles substantially block transmission of visible light through the display) and operate in a reflective mode, many electrophoretic displays can be made to operate in a so-called “shutter mode” in which one display state is substantially opaque and one is light-transmissive. See, for example, the aforementioned U.S. Pat. Nos. 6,130,774 and 6,172,798, and U.S. Pat. Nos. 5,872,552; 6,144,361; 6,271,823; 6,225,971; and 6,184,856. Dielectrophoretic displays, which are similar to electrophoretic displays but rely upon variations in electric field strength, can operate in a similar mode; see U.S. Pat. No. 4,418,346.\nAn encapsulated electrophoretic display typically does not suffer from the clustering and settling failure mode of traditional electrophoretic devices and provides further advantages, such as the ability to print or coat the display on a wide variety of flexible and rigid substrates. (Use of the word “printing” is intended to include all forms of printing and coating, including, but without limitation: pre-metered coatings such as patch die coating, slot or extrusion coating, slide or cascade coating, curtain coating; roll coating such as knife over roll coating, forward and reverse roll coating; gravure coating; dip coating; spray coating; meniscus coating; spin coating; brush coating; air knife coating; silk screen printing processes; electrostatic printing processes; thermal printing processes; ink jet printing processes; electrophoretic deposition (see US Patent Publication No. 2004/0226820); and other similar techniques.) Thus, the resulting display can be flexible. Further, because the display medium can be printed (using a variety of methods), the display itself can be made inexpensively.\nThe requirements for particles used in electrophoretic media are complex. The particles need to provide high optical density even when present only in very thin layers; because the switching speed of an electrophoretic medium (i.e., the time taken at a given voltage for the medium to change from one of its extreme optical states to the other) is proportional to the thickness of the medium, it is desirable to use a thin layer of the electrophoretic material itself, typically of the order of 50 μm or less, and in such thin media there may be only the equivalent of about a 2 μm thick layer of electrophoretic particles present. The electrophoretic particles need to maintain a stable charge, so that they will respond in a predictable manner to an electric field, and must not stick to surfaces with which they come into contact, nor (in the case of electrophoretic media containing particles with charges of both polarities) stick too firmly to particles of opposite charge. It may be difficult to find a single material which meets all the requirements for electrophoretic particles, and it is known to use as electrophoretic particles composite particles having a core of one material and a coating or shell of a different material. For example, the aforementioned U.S. Pat. Nos. 6,822,782, 7,002,728, and 7,230,750 describe electrophoretic particles having a core of a pigment, for example titania, carbon black or copper chromite, and a polymer coating. Other types of composite electrophoretic particles are described in U.S. Patent Application Publication No. 2002/0044333, and in Japanese Published Applications Nos. 2004-318,006; 2004-233831; 2004-029699; 2004-004469; 2000-227612; and 01-114829.\nAs already noted, it is desirable that electrophoretic particles provide high optical densities even when present only in very thin layers. To meet this requirement, electrophoretic particles are usually based on pigments; other types of materials, for example polymers colored with dyes, have been used but typically do not provide as high an optical density as is desirable. In electrophoretic media which require colored (as opposed to black and white) particles, it is often difficult to find pigments having a particular desired color. Inorganic pigments, while durable and well suited to use in electrophoretic media, offer a very limited color gamut. Organic pigments are available in a much wider range of colors, but frequently have a particle size much smaller than the wavelength of visible light, and so are of limited utility in light scattering media, such as electrophoretic media.\nAccordingly, it is desirable to find a way to permit small particle organic pigments to be used in electrophoretic media and displays, and the present invention provides a composite electrophoretic particle which enables this to be done."}
{"text": "Digital distribution of multimedia content (audio, video, etc.) and the impending convergence of industries that seek to make this goal a reality (computer, telecommunications, media, electric power, etc.) collide with the simplicity of making perfect digital copies. There exists a vacuum in which content creators resist shifts to full digital distribution systems for their digitized works, due to the lack of a means to protect the copyrights of these works. In order to make such copyright protection possible, there must exist a mechanism to differentiate between a master and any of its derivative copies. The advent of digital watermarks makes such differentiation possible. With differentiation, assigning responsibility for copies as they are distributed can assist in the support and protection of underlying copyrights and other “neighboring rights,” as well as, the implementation of secure metering, marketing, and other as yet still undecided applications. Schemes that promote encryption, cryptographic containers, closed systems, and the like attempt to shift control of copyrights from their owners to third parties, requiring escrow of masters and payment for analysis of suspect, pirated copies. A frame-based, master-independent, multi-channel watermark system is disclosed in U.S. patent application Ser. No. 08/489,172 filed on Jun. 7, 1995 and entitled “STEGANOGRAPHIC METHOD AND DEVICE” (now U.S. Pat. No. 5,613,004 issued Mar. 18, 1997), U.S. patent application Ser. No. 08/587,944 filed on Jan. 17, 1996 and entitled “METHOD FOR HUMAN-ASSISTED RANDOM KEY GENERATION AND APPLICATION FOR DIGITAL WATERMARK SYSTEM” (now U.S. Pat. No. 5,822,432 issued Oct. 13, 1998), and U.S. patent application Ser. No. 08/587,943 filed on Jan. 16, 1996 and entitled “METHOD FOR STEGA-CIPHER PROTECTION OF COMPUTER CODE” (now U.S. Pat. No. 5,745,569 issued Apr. 28, 1998). These applications describe methods by which copyright holders can watermark and maintain control over their own content. Any suspect copies carry all necessary copyright or other “rights” information within the digitized signal and possession of an authorized “key” and the software (or even hardware) described in these applications would make determination of ownership or other important issues a simple operation for the rights holder or enforcer.\nOptimizing watermark insertion into a given signal is further described in the U.S. patent application Ser. No. 08/677,435 filed on Jul. 2, 1996 and entitled “OPTIMIZATION METHODS FOR THE INSERTION, PROJECTION AND DETECTION OF DIGITAL WATERMARKS IN DIGITIZED DATA”, (now U.S. Pat. No. 5,889,868 issued Mar. 30, 1999). This application discloses accounting for the wide range of digitally-sampled signals including audio, video, and derivations thereof that may constitute a “multimedia” signal. The optimization techniques described in that application take into account the two components of all digitization systems: error coding and digital filters. The premise is to provide a better framework or definition of the actual “aesthetic” that comprises the signal being reproduced, whether through commercial standards of output (NTSC, CD-quality audio, etc.) or lossless and lossy compression (MPEG-2, Perceptual Audio Coding, AC-3, Linear Adaptive Coding, and the like), so that a watermark may be targeted at precisely the part of the signal comprising such an “aesthetic” in order that it be as robust as possible (i.e., difficult to remove without damaging the perceptual quality of the signal). However the content is stored, the signal still carries the digital watermark. Additionally, transmission media may be characterized as a set of “filters” that may be pre-analyzed to determine the best “areas” of the signal in which watermarks “should” be encoded, to preserve watermarks in derivative copies and ensure maximum destruction of the main, carrier signal when attempts are made to erase or alter the watermarked content.\nOptimal planning of digital watermark insertion can be based on the inversion of digital filters to establish or map areas comprising a given content signal's “insertion envelope.” That is, the results of the filter operation are considered in order to “back out” a solution. In the context of this discussion, the phrase “inverting” a filter may mean, alternatively, mathematical inversion, or the normal computation of the filter to observe what its effect would be, were that filter applied at a later time. Planning operations will vary for given digitized content: audio, video, multimedia, etc. Planning will also vary depending on where a given “watermarker” is in the distribution chain and what particular information needs that user has in encoding a given set of information fields into the underlying content. The disclosures described take into account discrete-time signal processing which can be accomplished with Fast Fourier Transforms that are well-known in the art of digital signal processing. Signal characteristics are also deemed important: a specific method for analysis of such characteristics and subsequent digital watermarking is disclosed in further detail in this application. The antecedents of the present invention cover time and frequency domain processing, which can be used to examine signal characteristics and make modifications to the signal. A third way would be to process with z-transforms that can establish signal characteristics in a very precise manner over discrete instances of time. In particular, z-transform calculations can be used to separate the deterministic, or readily predictable, components of a signal from the non-deterministic (unpredictable or random) components. It should be apparent to those skilled in the art that non-deterministic is a subjective term whose interpretation is implicitly affected by processing power, memory, and time restrictions. With unlimited DSP (digital signal processing) power, memory, and time to process, we might theoretically predict every component of a signal. However, practicality imposes limitations. The results of the z-transform calculations will yield an estimator of the signal in the form of a deterministic approximation. The difference between a signal reconstituted from the deterministic estimator and the real signal can be referred to as error, and the error in an estimator can be further analyzed for statistical characteristics. Those skilled in the art will be aware that Linear Predictive Coding (LPC) techniques make use of these properties. So the error can be modeled, but is difficult to reproduce exactly from compressed representations. In essence, this error represents the randomness in a signal which is hard to compress or reproduce, but in fact may contribute significantly to the gestalt perception of the signal.\nThe more elements of error determined with z-transforms, the better able a party is at determining just what parts of a given carrier signal are deterministic, and thus predictable, and what elements are random. The less predictable the watermark-bearing portion of a signal is and the more it contributes to the perception of the signal, as previously disclosed, the more secure a digital watermark can be made. Z-transform analysis would disclose just which phase components are deterministic and which are random. This is because it is difficult to compress or otherwise remove unpredictable signal components. Error analysis further describes the existence of error function components and would reliably predict what signals or data may later be removed by additional z-transform analysis or other compression techniques. In effect, the error analysis indicates how good an approximation can be made, another way of stating how predictable a signal is, and by implication, how much randomness it contains. Z-transforms are thus a specialized means to optimize watermark insertion and maximize the resulting security of encoded data from attempts at tampering. The results of a Z-transform of input samples could be analyzed to see “exactly” how they approximate the signal, and how much room there is for encoding watermarks in a manner that they will not be removed by compression techniques which preserve a high degree of reproduction quality.\nTime is typically described as a single independent variable in signal processing operations but in many cases operations can be generalized to multidimensional or multichannel signals. Analog signals are defined continuously over time, while digital signals are sampled at discrete time intervals to provide a relatively compact function, suitable for storage on a CD, for instance, defined only at regularly demarcated intervals of time. The accrued variables over time provide a discrete-time signal that is an approximation of the actual non-discrete analog signal. This discreteness is the basis of a digital signal. If time is unbounded and the signal comprises all possible values, a continuous-valued signal results. The method for converting a continuous-valued signal into a discrete time value is known as sampling. Sampling requires quantization and quantization implies error. Quantization and sampling are thus an approximation process.\nDiscreteness is typically established in order to perform digital signal processing. The issue of deterministic versus random signals is based on the ability to mathematically predict output values of a signal function at a specific time given a certain number of previous outputs of the function. These predictions are the basis of functions that can replicate a given signal for reproduction purposes. When such predictions are mathematically too complicated or are not reasonably accurate, statistical techniques may be used to describe the probabilistic characteristics of the signal. In many real world applications, however, determinations of whether a signal, or part of a signal, is indeed random or not is difficult at best. The watermark systems described in earlier disclosures mentioned above have a basis in analyzing signals so that analysis of discrete time frames can be made to insert information into the signal being watermarked. When signal characteristics are measured, a key factor in securely encoding digital watermarks is the ability to encode data into a carrier signal in a way that mimics randomness or pseudo randomness so that unauthorized attempts at erasing the watermark necessarily require damage to the content signal. Any randomness that exists as a part of the signal, however, should be estimated in order that a party seeking to optimally watermark the input signal can determine the best location for watermark information and to make any subsequent analysis to determine the location of said watermarks more difficult. Again, typical implementations of signal processing that use z-transforms seek to describe what parts of the signal are deterministic so that they may be described as a compact, predictable function so that the signal maybe faithfully reproduced. This is the basis for so-called linear predictive coding techniques used for compression. The present invention is concerned with descriptions of the signal to better define just what parts of the signal are random so that digital watermarks may be inserted in a manner that would make them more or less tamperproof without damage to the carrier signal. Additional goals of the system are dynamic analysis of a signal at discrete time intervals so that watermarks may be dynamically adjusted to the needs of users in such instances as on-the-fly encoding of watermarks or distribution via transmission media (telephone, cable, electric powerlines, wireless, etc.)\nSignal characteristics, if they can be reasonably defined, are also important clues as to what portion or portions of a given signal comprise the “aesthetically valuable” output signal commonly known as music or video. As such, perceptual coding or linear predictive coding is a means to accurately reproduce a signal, with significant compression, in a manner that perfectly replicates the original signal (lossless compression) or nearly replicates the signal (lossy compression). One tool to make better evaluations of the underlying signal includes the class of linear time-invariant (LTI) systems. As pointed out in Digital Signal Processing (Principles, Algorithms, and Applications), 3rd Ed. (Proakis and Manolakis), (also Practical DSP Modeling, Techniques, and Programming in C by Don Morgan) the z-transform makes possible analysis of a continuous-time signal in the same manner as discrete-time signals because of the relationship between “the convolution of two time domain signals is equivalent to multiplication of their corresponding z-transforms.” It should be clear that characterization and analysis of LTI systems is useful in digital signal processing; meaning DSP can use a z-transform and invert the z-transform to deterministically summarize and recreate a signal's time domain representation. Z-transforms can thus be used as a mathematical way in which to describe a signal's time domain representation where that signal may not be readily processed by means of a Fourier transform. A goal of the present invention is to use such analysis so as to describe optimal locations for watermarks in signals which typically have components both of deterministic and non-deterministic (predictable and unpredictable, respectively) nature. Such insertion would inherently benefit a system seeking to insert digital watermarks, that contain sensitive information such as copyrights, distribution agreements, marketing information, bandwidth rights, more general “neighboring rights,” and the like, in locations in the signal which are not easily accessible to unauthorized parties and which cannot be removed without damaging the signal. Such a technique for determining watermark location will help ensure “pirates” must damage the content in attempts at removal, the price paid without a legitimate “key.”\nSome discussion of proposed systems for a frequency-based encoding of “digital watermarks” is necessary to differentiate the antecedents of the present invention which processes signals frame-by-frame and may insert information into frequencies without requiring the resulting watermark to be continuous throughout the entire clip of the signal. U.S. Pat. No. 5,319,735 to Preuss et al. discusses a spread spectrum method that would allow for jamming via overencoding of a “watermarked” frequency range and is severely limited in the amount of data that can be encoded—4.3 8-bit symbols per second. Randomization attacks will not result in audible artifacts in the carrier signal, or degradation of the content as the information signal is subaudible due to frequency masking. Decoding can be broken by a slight change in the playback speed. It is important to note the difference in application between spread spectrum in military field use for protection of real-time radio signals versus encoding information into static audio files. In the protection of real-time communications, spread spectrum has anti jam features since information is sent over several channels at once, and in order to jam the signal, you have to jam all channels, including your own. In a static audio file, however, an attacker has all the time and processing power in the world to randomize each sub-channel in the signaling band with no penalty to themselves, so the anti jam features of spread spectrum do not extend to this domain if the encoding is sub-audible. Choosing where to encode in a super-audible range of the frequency, as is possible with the present invention's antecedents, can better be accomplished by computing the z-transforms of the underlying content signal, in order to ascertain the suitability of particular locations in the signal for watermark information.\nInstead of putting a single subaudible, digital signature in a sub-band as is further proposed by such entities as NEC, IBM, Digimarc, and MIT Media Lab, the antecedent inventions' improvement is its emphasis on frame-based encoding that can result in the decoding of watermarks from clips of the original full signal (10 seconds, say, of a 3 minute song). With signatures described in MIT's PixelTag or Digimarc/NEC proposals, clipping of the “carrier signal” (presently only based on results from tests on images, not video or audio signals which have time domains), results in clipping of the underlying watermark. Additionally, the present invention improves on previous implementations by providing an alternative computational medium to time/amplitude or frequency/energy domain (Fourier Transform) calculations and providing an additional measure by which to distinguish parts of a signal which are better suited to preserve watermarks through various DSP operations and force damage when attempts at erasure of the watermarks are undertaken. Further, the necessity of archiving or putting in escrow a master copy for comparison with suspect derivative copies would be unnecessary with the present invention and its proposed antecedents. Further, statistical techniques, not mathematical formulas, that are used to determine a “match” of a clip of a carrier signal to the original signal, both uneconomical and unreasonable, would not be necessary to establish ownership or other information about the suspect clip. Even if such techniques or stochastic processes are used, as in an audio spread-spectrum-based watermarking system being proposed by Thorn-EMI's CRL, called ICE, the further inability to decode a text file or other similar file that has been encoded using a watermark system as previously disclosed by above-mentioned U.S. patent applications including “Steganographic Method and Device”, “Method for Human-Assisted Random Key Generation and Application for Digital Watermark System”, “Method for Stega-cipher Protection of Computer Code”, and “Optimal Methods for the insertion, Protection and Detection of Digital Watermarks in Digitized Data”, where all “watermark information” resides in the derivative copy of a carrier signal and its clips (if there has been clipping), would seem archaic and fail to suit the needs of artists, content creators, broadcasters, distributors, and their agents. Indeed, reports are that decoding untampered watermarks with ICE in an audio file experience “statistical” error rates as high as 40%. This is a poor form of “authentication” and fails to establish more clearly “rights” or ownership over a given derivative copy. Human listening tests would appear a better means of authentication versus such “probabalistic determination”. This would be especially true if such systems contain no provision to prevent purely random false-positive results, as is probable, with “spread spectrum” or similar “embedded signaling”—type “watermarks,” or actually, with a better definition, frequency-based, digital signatures."}
{"text": "In a traditional Software as a Service (Saas) environment, each customer has its own copy of software, such as human resources (HR) software for each HR department of different corporate customers. By providing each corporate customer its own HR software copy, sensitive customer data accessed by the HR software is not shared by the corporate customers. However, each software copy consumes computing system resources and system administration resources. If there are thousands of customers, each software copy consumes many more computing system and system administration resources. Scaling to millions of customers further exacerbates the total cost in maintaining the SaaS environment. Accordingly, to alleviate such costs, a SaaS environment should provide a single copy of the software to multiple customers while also preventing sensitive data sharing between the multiple customers."}
{"text": "Electronics industry and its associated industries have known a rapid, spectacular development as time evolves. Further, all kinds of high-technology products and communication products related to computer technology are commercially available. The fast development of the new products not only advantageously influences our daily life and work, but also brings a great convenience to people. In response to all kinds of new information products being developed, especially notebook computers, most users become more critical with respect to the convenience in using them. Thus, whether electronic products manufactured in the future can provide a more convenient and effective service to customers will be an indicator to decide whether sale revenue and manufacturing technology of one large electronic company are higher and more advanced than other competitive ones respectively.\nFor meeting the increasing consumer demand about computers, a high-technology company has to not only maintain good quality of its products, but also increase the manufacturing speed. In brief, each portion of an assembly line, whether important or not, will affect the manufacturing speed. For example, the speed of installing a CD-ROM drive in a computer is one of important factors that affect the increase of the computer assembly speed. It is understood that competition among computer manufacturing companies is very fierce. Thus, one computer manufacturing company cannot have a large share of the computer market if both the computer assembly speed of the company is low and quality of computers produced by the company is poor. To the worse, the company may be eliminated from the competitive computer market.\nHowever, typically, a plurality of screws are employed to fasten a CD-ROM drive in the computer in the installation process. This has a drawback of low in assembly speed. Similarly, in a case of removing a malfunctioned CD-ROM drive for maintenance or replacement, a user has to sequentially detach the computer case, unfasten the screws by means of a screw driver, disconnect all power and data cables from the CD-ROM drive, and finally detach the CD-ROM drive from the computer case. As such, it is a tedious, time consuming, and labor intensive process during computer assembly. Similar drawbacks are also found in maintenance. This is an undesirable design. Further, it contradicts the trend of low cost, high efficiency, and mass production as pursued by modern industries.\nThus, it is desirable among vast consumers and computer manufacturers to provide a novel fastening mechanism capable of rapidly detaching a CD-ROM drive from a computer and installing the same in the computer in order to overcome the above drawbacks of the prior art."}
{"text": "1. Field of the Invention\nThe invention relates to devices and methods for treating incontinence. More specifically, particular embodiments of the invention use a variable control vacuum and/or vibratory device or other device to apply suction or other force to the male or female genital region, e.g. the region of the female urethral opening, the clitoral region and/or the vaginal region, or the male penile region.\n2. Description of Related Art\nA. Incontinence\nUrinary incontinence is a significant clinical problem and a major source of disability and dependency. Although urinary incontinence is particularly prevalent in the elderly, it affects all age groups. Millions of male and female adult Americans suffer from urinary incontinence, including at least one-half of all nursing home residents. The monetary costs of managing urinary incontinence, and the psychological costs associated with the problem, are great. Fecal incontinence presents similar concerns.\nThe most frequently occurring types of urinary incontinence are stress incontinence, urge incontinence, overflow incontinence, and mixed incontinence. Stress incontinence is a common form of incontinence in women. Intra-abdominal pressure exceeds urethral pressure upon coughing, sneezing, laughing, lifting, or like activity, causing leakage of urine. Physical changes associated with pregnancy, childbirth, and menopause, for example, are known to cause stress incontinence.\nUrge incontinence occurs when a patient loses urine while suddenly feeling the urge to urinate. The patient is unable to inhibit the flow of urine long enough to reach the toilet. Inappropriate bladder contractions are the most common cause of urge incontinence, and may occur in connection with central nervous system lesions, urinary infection, or bladder tumors, to name several examples.\nOverflow incontinence occurs when the bladder is unable to empty normally. Weak bladder muscles, caused e.g. by nerve damage from diabetes, or a blocked urethra, caused e.g. by tumors or urinary stones, are among the more common causes of overflow incontinence. Frequency or urgency involves the need or urge to urinate on an excessively frequent or habitual basis. Combinations of these and other types of incontinence, e.g. stress incontinence and urge incontinence, are often called mixed incontinence.\nMany options are available to treat incontinence in its various forms, including Kegel exercises, electrical stimulation, biofeedback, timed voiding or bladder training, medications, pessaries, implants, invasive or minimally invasive surgery, catheterization, and other methods and devices. Kegel exercises, or pelvic floor muscle exercises, for example, are intended to strengthen the muscles supporting the urethra, bladder, uterus and rectum. Weak pelvic floor muscles can contribute to urinary incontinence problems; hence strengthening those muscles can tend to alleviate those problems.\nAccording to several known surgical treatment methods, one or more stimulation systems, including electrodes controlled by an electronic control system, are implanted to electrically stimulate nerves controlling external sphincter and bladder functions. The electronic control system directs the administration of electric pulses to stimulate the nerves appropriately. Note U.S. Pat. No. 4,771,779, for example, which is incorporated herein by reference. U.S. Pat. No. 6,055,456, also incorporated herein by reference, shows an implantable medical lead for stimulation of the sacral nerves to treat incontinence.\nAlthough treatments requiring surgical intervention may be the preferred and most effective treatment mode in some situations, surgical intervention may be too extreme a measure in other situations. In some cases, surgical procedures to treat incontinence actually have a relatively low success rate; in many cases such procedures are irreversible. Additionally, a patient may hesitate to proceed with a surgical option, and/or a patient's physical condition may make surgical intervention inappropriate. Surgery may be inappropriate for pregnant patients, for example, or those of advanced age. Similarly, pharmacological treatment options may cause undesirable side effects and/or interactions with other medications. Non-surgical treatments, for example exercises or bladder training, may demand too high a degree of patient compliance or effort and thus may be resisted or otherwise ineffective.\nA need exists, therefore, to treat urinary and/or fecal incontinence in a non-invasive, non-pharmacological manner that is less likely to meet with patient resistance or cause physical trauma or side effects.\nB. Urogenital Anatomy\nThe female urethra passes just anterior to the vagina, and the external female urethral orifice is just posterior to the clitoris. The external urethral sphincter and the compressor urethrae compress the urethra and serve as functional sphincters. The pudendal nerve, which is derived from the anterior divisions of the ventral rami of S2 through S4, is the chief nerve of the perineum, providing sensation to the genitalia including the clitoris and controlling the motor function of the external urethral sphincter and the external anal sphincter. In the male, the pudendal nerve provides sensation from the penis and scrotum and provides motor control to the pelvic floor, including the anal sphincter.\nC. Clitoral Anatomy\nThe clitoris in the human female consists of a cylindrical, erectile organ composed of three parts: the outermost glans or head, the middle corpus or body, and the innermost crura. The glans of the clitoris is visualized as it emerges from the labia minora, which bifurcates to form the upper prepuce anteriorly and the lower frenulum posteriorly. The body of the clitoris consists of two paired corpora cavernosa of about 2.5 cm in length. The body extends under the skin at the corona to the crura. The two crura of the clitoris, formed from the separation of the most proximal portions of the corpora in the perineum, attach bilaterally to the undersurface of the symphysis pubis at the ischiopubic rami.\nA fibrous tunica albuginea ensheathes each corporal body made up of lacunar space sinusoids surrounded by trabecula of the vascular smooth muscle and collagen connective tissue. No retractor clitoridis muscle exists in humans as it does in other animals such as cattle and sheep, however a supporting suspensory ligament does hold the clitoris in the introital region.\nThe main arterial supply to the clitoris is from the ilio-hypogastric-pudendal arterial bed. The internal pudendal artery is the last anterior branch off the internal iliac artery. Distally, the internal pudendal artery traverses Alcock's canal, then terminates as it supplies the inferior rectal and perineal artery which supply the labia. The common clitoral artery continues to the clitoris. This artery bifurcates into a dorsal clitoral artery and a cavernosal clitoral artery.\nIn the normal female, autonomic efferent innervation of the clitoris passes from the pelvic and hypogastric nerves to the clitoris. Pelvic nerve stimulation results in clitoral smooth muscle relaxation and arterial smooth muscle dilation, causing an increase in clitoral cavernosal artery inflow and an increase in clitoral intracavernous pressure, which lead to tumescence and extrusion of the glans clitoris.\nThe clitoris has a dense collection of Pacinian corpuscles, Meissner's corpuscles, Merckel tactile disks, and free nerve endings. These sensory afferent nerves pass through the dorsal clitoral nerve to the pudendal nerve and into the sacral nerve roots (S2, S3, S4). (The male anatomy is completely homologous.) The presence of at least one somatic reflex arc is demonstrated by the bulbo-cavernosal reflex; squeezing the clitoris causes the anal sphincter to contract—the so-called “anal wink.”\nD. Female Sexual Dysfunction\nClitoral erectile insufficiency or reduced clitoral arterial flow may be caused by atherosclerosis, diabetes, or age-related causes, among other factors. Reduced clitoral arterial flow may lead to fibrosis of the clitoral cavernosa and reduced clitoral physiological function. In an animal model, Park et al. demonstrated that significant collagen synthesis occurs when the arterial inflow to the clitoris is compromised. This work demonstrated the importance of maintaining arterial flow to the clitoris to prevent collagen synthesis and fibrosis on the smooth muscle. See Park, K., et al., Vasculogenic Female Sexual Dysfunction: The Hemodynamic Basis for Vaginal Engorgement Insufficiency and Clitoral Erectile Insufficiency, IJIR, 9:27-37, 1997.\nIt is believed that the difficulty or inability to achieve clitoral tumescence may be related to and associated with other symptoms of female sexual arousal disorder. According to the International Consensus Report on Female Sexual Dysfunction, Female Sexual Arousal Disorder (FSAD) is defined as the persistent or recurrent inability to attain or maintain adequate genital lubrication or swelling responses resulting in personal distress. FSAD may be expressed as a lack of subjective excitement or lack of genital (lubrication/swelling) or other somatic responses (AFUD Consensus Report of FSD, 1998).\nA non-pharmacological approach to treatment that causes blood flow and engorgement, thereby applying a stimulus to the sensory nerve endings in the clitoris, periurethral area, and/or genital area, would be very beneficial to a large group of women complaining of FSAD."}
{"text": "Although many of recent mobile phones are integrally equipped with a camera, it is difficult with the employment of the camera installed in such a mobile phone to take a clear picture of image especially when a subject exists in a dark place because of the performance of the camera. Namely, although the image pickup performance of the camera depends largely on brightness as well as color rendering of flash illumination to be emitted from a flash mounted on the mobile phone, no one has succeeded as yet to develop such a flash which is capable of exhibiting a sufficient brightness and satisfactory color rendering for enabling a clear picture of image to be taken under such a dark condition.\nIn recent years, in view of the possibility of reducing the size and weight of a portable terminal unit, much attention has been paid to the employment of an LED flash. The LED flash to be mounted on a portable terminal unit is required, first of all, to be excellent in brightness. Because of this, a plurality of LED chips are generally employed for the purpose of flash illumination at present. However, since the capacity of battery to be employed therein is limited, the efficiency of each LED chip is desired to be as high as possible. Furthermore, since the direction of light toward a subject to be photographed, i.e., the directivity of brightness is a key if the LED chip is to be employed for the purpose of flash illumination, the luminescence intensity (cd) of the LED chip is one of the important characteristics thereof.\nTherefore, there have been proposed various ideas such, for example, as; an idea to dispose a reflection plate close to and around an LED chip for the purpose of enhancing the luminescence intensity thereof; an idea to dispose a plurality of LED chips so as to circularly surround one LED chip (see, for example, JP-A 2005-158958); and an idea to dispose a plurality of LED chips with centers thereof being linearly arrayed (see, for example, JP-A 2005-109212 and JP-A 10-22529). In spite of these proposals, no one has succeeded as yet to develop such a flash which is capable of exhibiting satisfactory brightness."}
{"text": "The present invention relates to a focusing unit for controlling the focusing voltage of a cathode ray tube within a TV receiver by appropriately controlling the high voltage of the FBT (flyback transformer) and, more particularly, to a dynamic focusing unit capable of controlling the peripheral focus in the up and down and the right and left dimensions of the TV screen correctly and over a wider range.\nGenerally, the focusing unit is a device for focusing the electron beam, being an electrostatic lens formed between anode electrode and cathode electrodes by applying pertinent voltage to the focusing electrodes such that a clear image is displayed on a fluorescent screen by the electron beam emitted from the electron gun of the picture tube.\nFIG. 8 is a schematic structural diagram of a picture tube with a spherical screen controlled by the focusing circuit of Korean Utility Model official publication gazette No. 89-8220, and wherein the electrons emitted from cathode 100 of the electron gun are made into an electron beam being radiated to a beam focusing point on a fluorescent screen 130. For focusing this electron beam, focusing electrode 110 is provided within the interior of the picture tube and then pertinent voltage is applied from the exterior. Thereby, the electron beam is focused by the electrostatic field of focusing electrode 110.\nHowever, although the electron beam focused upon the fluorescent screen 130 of the aforementioned spherical picture tube does not produce deviation, in the case when the fluorescent screen 130 is flat as the picture screen 140, a focusing deviation 150 is produced at the peripheral regions of the flat screen.\nAccordingly, since the conventional focusing device was a device designed for a spherical cathode ray tube, a focusing deviation has been produced at the peripheral region of the flat cathode ray tube.\nIn order to solve the above-described defect, a device has been developed in which focusing is correctly controlled at peripheral regions of cathode ray tubes by utilizing the horizontal flyback pulse signal and vertical output signal.\nAs shown in FIG. 9, a focusing device for controlling the electron beam to be focused on the flat screen of cathode ray tube 300 in response to the signal applied to focusing electrode 110 comprises a vertical signal integrating circuit 10 for integrating vertical pulse signal S.sub.1, thereby converting to a vertical parabolic waveform signal S.sub.2, a vertical amplifying circuit 20 for variably amplifying the output signal of said vertical signal integrating circuit 10, thereby controlling the up and down focusing on the screen, a horizontal signal integrating circuit 30 for integrating the horizontal pulse signal S.sub.4 inputted after being extracted from the flyback transformer, thereby converting to a horizontal parabolic waveform signal S.sub.5, a buffer circuit 40 for variably current amplifying the output signal of said horizontal signal integrating circuit, thereby controlling the right and left focusing on the screen, and an inversion amplifying circuit 50 for inversely amplifying the mixed vertical and horizontal parabolic waveform signals S.sub.6. The signal S.sub.7 outputted from said inversion amplifying circuit 50 is applied to focusing electrode 110 of picture tube 300 and also to anode 160 of picture tube 300. These elements are arranged so that the phenomenon of deviating focus at peripheral regions of the screen is prevented. The focusing unit serves a function which controls appropriately the high voltage coming out of the flyback transformer (FBT) thereby controlling the focus of the cathode ray tube within the TV receiving set.\nHowever, there are many kinds of cathode ray tubes and different focusing voltages are required by each kind. In order to feed focusing voltage of the cathode ray tube for each kind, the wider the variable range of the voltage of the focusing unit is, the better.\nA conventional focusing unit is shown in FIGS. 6 and 7 and comprises a knob 208 passing through hole 201 in the rectangular main body 200 of the focusing unit having annular ring 204 protruded with two rotation limiting protrusions 202, 203 extending inward, contact terminal fixing means 206 having a horseshoe shape to receive variable contact pin 205, knob body 210 structured with variable shaft means 207 formed integrally with said fixing means 206 as well as knob 208, and ceramic base board 230 provided with input and output voltage terminals 231, 232, circular printed portion 233 and bent printed portion 234.\nHowever, in this conventional focusing unit, since variable contact pin 205 is moved along circular printed portion 233 around knob 208, the extent of the output voltage terminal 231 from the input voltage 25 terminal 232 is limited within a narrow range.\nTherefore, since it has a simple structure that the first focusing voltage variable range 240 cannot cover by less than 10%-15% and also the second focusing voltage variable range 250 cannot cover by less than 2%-3%, there has been a problem of achieving a satisfactory variable voltage focusing range."}
{"text": "Many cities and towns have thousands of streetlights to keep streets and walkways lit at night. A given streetlight may be turned on or off using a photocell deployed on the streetlight."}
{"text": "The purpose of a catalytic converter for an internal combustion engine or a gas turbine is to convert pollutant materials in the exhaust, e.g., carbon monoxide, unburned hydrocarbons, nitrogen oxide, etc., to carbon dioxide, nitrogen and water. Conventional catalytic converters utilize a ceramic honeycomb monolith having square or triangular, straight-through openings or cells with catalyst deposited on the walls of the cells; catalyst coated refractory metal oxide beads or pellets, e.g., alumina beads; or a corrugated thin metal foil monolith, e.g., a ferritic stainless steel foil or a nickel alloy foil, having a catalyst carried on or supported on the surface. The catalyst is normally a noble metal, e.g., platinum, palladium, rhodium, ruthenium, or a mixture of two or more of such metals, and/or zeolite coatings. The catalyst catalyzes a chemical reaction, mainly oxidation, whereby the pollutant is converted to a harmless by-product which then passes through the exhaust system to the atmosphere.\nHowever, conversion to such harmless by-products is not efficient initially when the exhaust gases are relatively cold. To be effective at a high conversion rate, the catalyst and the surface of the converter with which the gases come in contact must be at or above a minimum temperature, e.g., 390.degree. F. for carbon monoxide, 570.degree. F. for volatile organic compounds (VOC) and 1000.degree. F. for methane or natural gas. Otherwise, conversion to harmless by-products is poor and cold start pollution of the atmosphere is high. Once the exhaust system has reached its normal operating temperature, the catalytic converter is optimally effective. Hence, it is necessary for the relatively cold exhaust gases to make contact with a hot catalyst so as to effect satisfactory conversion. Compression ignited engines, spark ignited engines and reactors in gas turbines have this need.\nTo achieve initial heating of the catalyst at or prior to engine start-up, there is conveniently provided an electrically heatable catalytic converter, preferably one formed of a thin metal monolith, either flat thin metal strips, straight corrugated thin metal strips, pattern corrugated thin metal strips, (e.g., herringbone or chevron corrugated) or variable pitch corrugated thin metal strips (See U.S. Pat. No. 4,810,588 dated Mar. 7, 1989 to Bullock et al), or a combination thereof, which monolith is connected to a voltage source, e.g., a 12 volt to 108 volt, or higher, power supply, preferably at the time of engine start-up and afterwards to elevate and maintain the catalyst to at least 650.degree. F. plus or minus 20.degree. F. Alternatively, power may also be supplied for 5 to 10 or so seconds prior to start-up of the engine. Catalytic converters containing a corrugated thin metal (stainless steel) monolith have been known since at least the early seventies. See Kitzner U.S. Pat. Nos. 3,768,982 and 3,770,389 each dated Oct. 30, 1973. More recently, corrugated thin metal monoliths have been disclosed in U.S. Pat. No. 4,711,009 dated Dec. 8, 1987; U.S. Pat. No. 4,381,590 to Nonnenmann et al dated May 3, 1983, copening application U.S. Ser. No. 606,130 filed Oct. 31, 1990 by William A. Whittenberger and entitled Electrically Heatable Catalytic Converter, now U.S. Pat. No. 5,070,694 dated Dec. 10, 1991, and commonly owned with the present application, and International PCT Publication Numbers WO 89/10470 and WO 89/10471 each filed Nov. 2, 1989. However, a common problem with many such prior devices has been their inability to survive severe automotive industry durability tests which are known as the Hot Shake Test and the Hot Cycling Test.\nThe Hot Shake Test involves oscillating (100 to 200 Hertz and 28 to 60 G inertial loading) the device in a vertical attitude at high temperature (between 800.degree. and 950.degree. C; 1472.degree. to 1742.degree. F., respectively) with exhaust gas from a running internal combustion engine simultaneously being passed through the device. If the electrically heatable catalytic device telescopes or displays separation or folding over of the leading, or upstream, edges of the foil leaves up to a predetermined time, e.g., 5 to 200 hours, the device is said to fail the test. Usually a device will fail within 5 hours if it is going to fail. Five hours is equivalent to 1.8 million cycles at 100 Hertz.\nThe Hot Cycling test is conducted on a running engine at 800.degree. to 950.degree. C. (1472.degree. to 1742.degree. F.) and cycled to 120.degree. to 150.degree. C. once every 15 to 20 minutes, for 300 hours. Telescoping or separation of the leading edges of the foil strips is considered a failure.\nThe Hot Shake Test and the Hot Cycling Test, hereinafter called \"Hot Tests\", have proved very difficult to survive, and many efforts to provide a successful device have been either too costly or ineffective for a variety of reasons.\nPreviously tested samples of EHC's in automotive service and comprised entirely of heater strips in electrical parallel, did not have adequate endurance in Hot Tests or did they satisfy the need for lower power ratings. In repeated efforts to arrive at a suitable design using purely parallel circuit construction, samples were made and tested with a wide range of parameters, including a length-to-diameter aspect ratio of from 0.5 to 1.5, cell densities of from 100 to 500 cells per square inch, individual strip heaters as long as 15 inches, and parallel circuits limited to as few as 2 heater strips.\nDevices made according to these design parameters proved unsatisfactory in the Hot Tests because (a) terminal resistance was too low and therefore, the devices drew too much power, (b) the relatively high voltage differential between laminations associated with small numbers of parallel heater strips caused some arcing and, (c) Hot Tests could not be passed consistently. With regard to (c), EHC's with heater strips longer than about 7\" measured from the center of the core to the outside of the core have generally not passed the Hot Shake Test. Resistance that is too low causes one or more of the following problems: (a) the weight and size of the battery become unacceptably high and/or expensive; (b) the EHC has to be made so large in diameter that longer heater strips have to be used which induces a tendency to fail the Hot Tests. The devices were, of course, useful in less stringent environments.\nPrior structures, such as that described in U.S. Pat. No. 4,928,485 have had all of the corrugated thin metal heater strip members connected in a manner such that all of the strips extended spirally outwardly from a central electrode to a circular shell which served as the electrode of opposite polarity. The strips served as heaters for the core. However, power levels of less than 2.0 kilowatts at 7 volts measured at the EHC terminals could not be achieved when all of the heater strips were in parallel because the terminal resistance of the EHC was too low.\nIn the following description, reference will be made to \"ferritic\" stainless steel. A suitable formulation for ferritic stainless steel alloy is described in U.S. Pat. No. 4,414,023 to Aggen dated Nov. 8, 1983. A specific ferritic stainless steel useful herein contains 20% chromium, 5% aluminum, and from 0.002% to 0.05% of at least one rare earth metal selected from cerium, lanthanum, neodymium, yttrium, and praseodymium, or a mixture of two or more of such rare earth metals, balance iron and trace steel making impurities. Another metal alloy especially useful herein is identified as Haynes 214 which is commercially available. This alloy is described in U.S. Pat. No. 4,671,931 dated Jun. 9, 1987 to Herchenroeder et al. Haynes 214 is characterized by high resistance to corrosion. A specific example contains 75% nickel, 16% chromium, 4.5% aluminum, 3% iron, trace amounts of one or more Rare Earth metals, 0.05% carbon, and steel making impurities. Ferritic stainless steel and Haynes 214 are examples of high temperature resistive, corrosion resistant metal alloys useful herein. Suitable alloys must be able to withstand temperatures of 900.degree. C. to 1100.degree. C. over prolonged periods.\nOther high temperature resistive, corrosion resistant metal alloys are known and may be used herein. The thickness of the metal foil heater strips should be in the range of from 0.001\" to 0.005 preferably 0.001\" to 0.002\". The term \"thin\" as used herein means a thickness of the metal strip in the range of from 0.0015\" to about 0.005\". The two alloys mentioned above have been found most satisfactory for use in the present invention, especially Haynes 214.\nIn the following description, reference will also be made to fibrous ceramic mat, woven ceramic fiber tape, or insulation. Reference may be had to U.S. Pat. No. 3,795,524 dated Mar. 5, 1974 to Sowman and to the Hatch U.S. Pat. No. 3,916,057 dated Oct. 28, 1975 for formulations and manufacture of ceramic fibers and mats useful herein. One such ceramic fiber material is currently available from 3-M under the registered trademark \"NEXTEL\" 312 Woven Tape and is especially useful herein. Ceramic fiber mat is available as \"INTERAM\" also from 3-M.\nIt is a principal object of the present invention to provide a core element which enables simplified, less costly manufacture of a core useful in a catalytic converter."}
{"text": "Linear voltage regulators are used in many electronic devices to maintain a constant voltage to drive a load. Because the load may vary, the resistance of the regulator varies based on the load so that a constant voltage is produced. One type of linear voltage regulator is a low-dropout (LDO) regulator. LDO regulators are configured to regulate voltage even if the supply voltage is close to the regulated output voltage. LDO regulators include a differential amplifier that drives a power transistor. One input of the differential amplifier is a reference voltage generated by a voltage reference (e.g., a bandgap reference, a Zener diode, etc.). The second input of the differential amplifier is a fraction of the output voltage of the LDO regulator. Thus, the drive voltage to the power transistor changes to regulate the output voltage based on the value of the output voltage itself. In other words, if the output voltage rises relative to the reference voltage, the drive voltage powering the power transistor also changes to maintain an output voltage that remains constant."}
{"text": "The present invention relates to novel pharmacological agents, particularly prostaglandin analogs. The present invention provides novel 2a,2b-dihomo-15-alkyl-PGF.sub.2.alpha. compounds, a description of which, including their preparation and pharmacological uses, is incorporated here by reference from U.S. Pat. No. 3,974,195."}
{"text": "There has been conventionally known a printer which can print, using an application program from a host even if the application program and a device driver are not installed in the host.\nHowever, if a user wants to execute the application program in a new specification, the application program needs to be updated."}
{"text": "PDM systems manage PLM and other data. Improved systems are desirable."}
{"text": "1. Technical Field\nThe present invention relates generally to processor performance measurement systems, and more particularly, to a performance counting system with low latency and error.\n2. Description of the Related Art\nPerformance measurement in processing units is typically performed by a set of counters that accumulate information about the usage of units within the processor. By obtaining such usage information, sophisticated power management algorithms can adjust operating conditions such as processor frequency and voltage, idle time or other energy usage control mechanism to reduce power consumption or dissipation, while ideally providing a minimal impact on processing performance.\nAccurate performance evaluating systems have been proposed based on a complex polynomial metric, but require measurement of a large number concurrent events. If the events contributing to the performance measurement are the individual usage of a set of processing units as mentioned above, then a counter is provided for each unit for which usage is being measured. In such implementations, a large number of counters is required to measure performance in present-day processors, which may have multiple pipelines having dozens of processing units, cache units, and other units that contribute to overall processor performance. However, existing performance counting implementations typically require a separate counter for each metric being measured and typically do not process a large number of events concurrently. Therefore, such performance counting implementations cannot provide an accurate estimate of performance change when power management control changes operating parameters of a processor.\nOne alternative to the above-described scheme of providing a usage counter for each metric, is to time-multiplex a smaller set of counters (or a single counter) to perform the individual measurements. However, the overall latency of such a scheme is too high for energy management systems requiring a fine granularity of control. Further, the inter-metric measurement delays introduce error into the measurements. Even when individual counters are employed, they are not typically accessed simultaneously, which also introduces both inter-metric delay error and latency.\nIt is therefore desirable to provide a performance counting method and system that provides usage information having low error and latency. It would be further desirable to provide such a system and method that does not require a counter for each metric contributing to the overall performance measurement."}
{"text": "Ion exchange resin beads are used as an anion or cation exchanger in many applications, such as water treatment, production of pure water, catalysts, and purification of sugar solutions. Ion exchange resins for use in these applications and copolymers to be converted ion exchange resins generally include a gel type resin having a non-porous structure and a macroporous type resin having a porous or macro-network structure. Superiority of the ion exchange resins of gel type to those of porous type resides in higher physical strength as represented by friability and higher exchange capacity. Superiority of the porous type ion exchange resins to those of gel type resides in higher physical stability as represented by osmotic pressure resistance and higher rate of ion-exchange reaction owing to their greater specific surface area in contact with a substance to be treated. Having a porous or macro-network structure, the macroporous type ion exchange resins are also useful for water treatment or removal of impurities from stock water for the production of pure water.\nHowever, structurally lacking in exchange capacity and physical strength as compared with gel type ion exchange resins, the macroporous type ion exchange resin beads are liable to destruction by various stresses exerted thereon under conditions of use.\nIn order to overcome the respective disadvantages of the conventional ion exchange resins, there has been proposed a two-stage polymerization process for producing an ion exchange resin possessing advantageous characteristics of both gel type and porous type resins or a copolymer capable of providing such an ion exchange resin, comprising swelling a crosslinked copolymer having a macro-network structure with a monofunctional monomer and a polyfunctional monomer by imbibing and then polymerizing the monomers, as disclosed, e.g., in JP-B-48-17988 (the term \"JP-B\" as used herein means an \"examined Japanese patent publication\") and U.S. Pat. No. 3,991,017. However, when a functional group is introduced into the thus produced crosslinked copolymer by known techniques, an ion exchange resin having as much strength as can reasonably be expected from the total degree of crosslinking cannot be obtained. This appears to be because the monofunctional monomer and the polyfunctional monomer imbibed and filled into pores of the crosslinked copolymer having a macro-network structure has a weak interaction with the crosslinked copolymer and also because the imbibing is non-uniform only to produce a non-uniform crosslinked copolymer.\nExamples of the water treatments with an ion exchange resins include removal of impurities from a condensate of cooling water for the core of a nuclear reactor or a condensate of the steam generated in the reactor. For example, in boiling water reactors (hereinafter abbreviated as \"BWR\") where highly purified water must be used as a primary coolant, several devices for purifying the spent water are provided in the course of the circulation equipment so that the water should be repeatedly used.\nFor better understanding, a basic construction of BWR is explained below by referring to FIG. 1. The steam generated in reactor 1 is forwarded to turbine 3 through main steam pipe 2. After working, the steam is transferred to condenser 4 where it is condensed. The condensate is then forwarded to condensate purifying apparatus 5 by pump 6. The thus purified water is fed to feedwater heater 8 by pump 7 where it is heated. Most of the heated water is supplied to reactor 1 through water pipe 9, while a part of the heated water is sent to condensate purifying apparatus 11 provided on reactor recirculation line 10.\nNuclear-power generation of the above-described type requires a large quantity of water having extremely high purity for reducing radioactivity induction. Criteria of water quality for nuclear-power generation are as follows.\n______________________________________ Item Criteria ______________________________________ Conductivity (at 25.degree. C.) &lt;0.1 .mu..OMEGA./cm Metallic impurities &lt;15 ppb Chloride ion &lt;1 ppb ______________________________________\nSince general tap water contains about 10 ppm of metals, it must be purified about 1000-fold so as to reduce the metals below 15 ppb. The condensate additionally contains trace amounts of metallic impurities released from the nuclear power plant. Such metallic impurities include ionic impurities and noinionic impurities called clad mainly composed of amorphous iron hydroxide and iron oxyhydroxide. These impurities are generally removed in a condensate purifying apparatus (e.g., as shown by numerals 5 and 11 in FIG. 1) using a demineralizer with an inlet and an outlet which is filled with a basic ion exchange resin and an acidic ion exchange resin.\nResins conventionally employed in such a condensate purifying apparatus are gel type ion exchange resins having a particle size of from about 0.1 to 1 mm and a specific surface area of about 0.03 m.sup.2 /g. The conventional ion exchange resins have a sufficient amount of ion exchange groups for removing ionic impurities in the condensate, and the thus purified water fulfills the above-specified criteria in most cases. However, nonionic impurities such as clad are difficult to remove with ion exchange resins, and some of the conventional gel type ion exchange resins attain a very low rate of removal of the nonionic impurities. Should the nonionic impurities remain in the condensate, they are activated in the reactor or adhered to the reactor to increase the radiation dose, and thus, removal of the nonionic impurities is of great importance. In particular, a fresh ion exchange resin has a small specific surface area and has therefore a small capacity of removing clad and particularly a low rate of removal. An old resin having been used for 3 to 4 years has an increased capacity of removal of clad probably due to swelling of the surface and so increased specific surface area.\nIt is known, on the other hand, that an ion exchange resin having its surface roughened through oxidation treatment or mechanical abrasion and thereby having an increased specific surface area is suitable for removal of clad as taught in JP-A-2-273550 and JP-A-2-307534. This is because clad, which mainly composed of amorphous iron hydroxide or iron oxyhydroxide and are therefore difficult to remove with an ion exchange resin, is easily dissolved in an acidic atmosphere created on the surface of an ion exchange resin to become an iron ion which is ready to be trapped on the surface of the ion exchange resin. Since the resin once catches an iron ion on the surface thereof, it is important for the resin to have an ability of entrapping the iron into the inside thereof. However, the above-described resin having its surface roughened by oxidation treatment is disadvantageous in that the strength is lessened by oxidation treatment and that the oxidation treatment is no more than accelerated deterioration, i.e., making the surface of the ion exchange resin particles substantially old like after repeated use. On the other hand, a porous type resin, although having a large specific surface area, has too small strength to withstand use in place of the gel type resin.\nU.S. Pat. No. 4,564,644 discloses a gel type resin having a core/shell structure, and U.S. Pat. No. 4,975,201 proposes to use this resin for removal of clad from a condensate. Such a resin includes, as illustrated in the working examples, a styrene-divinylbenzene copolymer skeleton having bonded thereto an ion exchange group, the skeleton resin being produced by imbibing a monomer mixture of styrene, divinylbenzene, and a free radical initiator into seed particles of a styrene-divinylbenzene copolymer, polymerizing from 40 to 95% of the absorbed monomers within the seed particles, and then continuously adding thereto a styrene monomer containing substantially no free radical initiator under such a condition that the styrene monomer be absorbed into the copolymer particles and polymerized therein. It is stated that such an ion exchange resin has high strength and stability against osmotic pressure and is therefore suited for use in treatment of a condensate from a BWR reactor. Nevertheless, production of the resin having such a core/shell structure in which the polymer composition gradually varies from the center to the outer surface involves about three different polymerization steps each requires proper control of the degree of polymerization. In particular, the third step of polymerization must be carried out by adding a monomer in small portions under a constant condition."}
{"text": "This invention relates to a position measuring apparatus for an underground excavator, which is useful in measuring a position of excavation by the underground excavator, such as a pipe-jacking machine or a shield machine, which is advancing in the ground while excavating a tunnel, to an optical deflection angle measuring apparatus for measuring with lights a deflection angle between two line segments connecting a cardinal point to respective points set on opposite sides of the cardinal point with distances left from the cardinal point, and also to a position measuring apparatus for an underground excavator, which by using this optical deflection angle measuring apparatus, measures a position of excavation by the underground excavator advancing along a curved path.\nFor an underground excavator such as a pipe-jacking machine or a shield machine which advances in the ground while excavating a tunnel, it is necessary to enable the excavator to accurately advance along a planned line which is a preset advancing route. For this purpose, it is desired to make it possible to measure the current position of the advancing underground excavator in real time and accurately. Namely, if the underground excavator begins to advance off the planned line, an operator can quickly find out this problem and can take a countermeasure at an early stage if the operator is provided in real time with highly-reliable information on the current position of the underground excavator. Control can therefore be performed with ease to make the underground excavator advance along the planned line, and improvements in the accuracy of the work can also be expected. As techniques for measuring a position of excavation by an underground excavator, a variety of methods have been used to date including a method in which xe2x80x9ca measurement is performed manually by using a transitxe2x80x9d, a method in which xe2x80x9ca position of excavation by an underground excavator is measured by arranging an oscillating coil, which generates an induction field, on the underground excavator and measuring the intensity of the induction field with a receiving coil arranged on the groundxe2x80x9d, and a method in which xe2x80x9ca position of excavation by an underground excavator is measured by laying a cable way conversely on the ground, causing a current to flow through the cable way to generate an induction field, and detecting the intensity of the induction field with a receiving coil arranged on the underground excavatorxe2x80x9d. However, these conventional techniques for the measurement of positions of underground excavators are inherently incapable of performing real-time measurements of positions of excavation or, even if such real-time measurements are feasible theoretically, can hardly be put into practical use.\nAs a position measuring technique for an underground excavator, which can reduce these problems, a direction detecting apparatus for a shield machine has been proposed as disclosed in JP 61-45092 A. This direction detecting apparatus for the shield machine, which has been proposed previously, is an apparatus in which xe2x80x9ca first laser beam oscillator for performing radiation into a forward tunnel and a first laser beam receiver capable of receiving a laser beam from a front are mounted on a base such that they can be turned in an X direction (yawing direction) and an Y direction (pitching direction) by servomotors, a measuring instrument capable of detecting their turned angles by sensors is arranged in an entrance section to the tunnel, a second laser beam oscillator for performing radiation toward a first laser beam receiver at a rear, a second laser beam receiver capable of receiving a laser beam from the first laser beam at the rear and a third laser beam oscillator for performing radiation toward the shield machine at a front are arranged on a base such that they can be turned in the X direction and the Y direction by servomotors, a measuring relay designed to detect their turned angles is arranged in an intermediate section of the tunnel, and a third laser beam receiver and a pitching-rolling meter, which can receive a laser beam from the second laser beam oscillator at the rear and can detect the X and Y directions and a rolling angle, respectively, are arranged on the shield machinexe2x80x9d.\nUpon performing construction work on the ground or under the ground, it is necessary to determine an angle relationship between two line segments which connect a point, which serves as a cardinal point, to points set on opposite sides of the cardinal point with distances left from the cardinal point. Upon constructing a curved or bent road, it is necessary to determine the angle of the curved or bent part of the road surface under construction. In this case, a measurement point which serves as a cardinal point is set at a suitable place in the curved or bent part, measurement points are also set in road-surface working zones on the opposite sides of the cardinal point with distances left from the cardinal point, and an angle between line segments connecting the measurement point as the cardinal point to the respective measurement points on the opposite sides of the cardinal point is measured. To excavate a curved tunnel by an underground excavator, the position of excavation must be ascertained to determine whether or not the underground excavator is accurately advancing along a planned route (preset advancing route). In this case, it is also necessary to determine an angle relationship between line segments connecting a measurement point, which serves as a cardinal point, to respective measurement points on opposite sides of the cardinal point, as will be described in detail subsequently herein. To determine an angle relationship between line segments on opposite sides of such a cardinal point as an apex, either an interior angle or an exterior angle of both the line segments can be measured, and this object can be achieved if a value concerning such an angle as enabling to automatically specify the angle relationship can be measured. In this specification, a value concerning such an angle as enabling to specify the angle relationship between such two line segments is called a xe2x80x9cdeflection anglexe2x80x9d.\nIn construction work, a measuring method making use of a transit has been adopted to date in general to measure such deflection angles. As this measuring method of a deflection angle by a transit is a method which relies upon human ability, it requires labor such as a skilled worker and moreover, takes a long time for each measurement. In addition, a turning mechanism is required to turn a telescope in a yawing direction (horizontal direction) and a pitching direction (vertical direction), so that mechanical measurement errors tend to occur due to this turning mechanism, thereby making it difficult to assure high measuring accuracy. Furthermore, if external force such as that causing tilting in the yawing direction and/or pitching direction, such as vibrations, is applied, a measurement error takes place by such external force, thereby affecting the results of the measurement.\nIncidentally, upon excavating a curved tunnel by an underground excavator which advances in the ground while excavating a tunnel, the position of excavation by the underground excavator is measured so that the underground excavator can accurately advance along a planned route. Examples of underground excavators of this type can include small-diameter pipe-jacking machine for burying in the ground small-diameter pipes which men cannot enter and semi-shield and shield machines for burying in the ground large-diameter pipes which men can enter. To measure the position of excavation by such an underground excavator, it is the general practice to set a measuring start point and a measuring end point at a positionxe2x80x94which serves as a start pint for excavation by the underground excavator, such as a starting shaftxe2x80x94and within the underground excavator, respectively, and also to set a suitable number of intermediate measuring points between the measuring start point and the measuring end point as the excavation by the underground excavator proceeds. By measuring distances between these measuring points and also measuring deflection angles between line segments connecting each intermediate measuring point and its adjacent measuring points on opposite sides of the intermediate measuring point, the position of excavation by the underground excavator is determined by computation on the basis of the results of these measurements. When measuring such deflection angles in the course of the measurement of positions of excavation by the underground excavator, the method in which measurements are performed by a transit has also been adopted to date. The method which makes use of a transit, however, takes a long time for each measurement and requires labor as mentioned above and, especially when measurements are performed in a narrow tunnel, the measurement work involves a great deal of labor and danger. As techniques for measuring the position of excavation by an underground excavator of this type, there are techniques which have adopted a method which optically measures a deflection angle by a laser beam without relying upon a transit.\nAs a representative example of the techniques which had adopted such a method for the measurement of the position of excavation by an underground excavator, the technique disclosed in JP 5-340186 A can be mentioned by way of example. According to the technique disclosed in JP 5-340186 A (hereinafter called xe2x80x9cthe conventional techniquexe2x80x9d), xe2x80x9ca laser aiming system having an angle measuring function is arranged in front of a rear collimation point set in a curved tunnel, a position detecting element (photoelectric element) provided with a miniature reflecting prism is arranged as a target within a shied machine, and a suitable number of wedge prismsxe2x80x94each of which is equipped with a range finder capable of measuring the direction change angle of a laser beam obtained by refracting a laser beam from the laser aiming systemxe2x80x94are arranged as the excavation by the underground excavator proceeds.\nUpon measuring the position of excavation by an underground excavator by this conventional technique, each wedge prism is caused to turn by a remote operation such that a laser beam from a laser aiming system is always directed through the wedge prism onto a target in the shield machine. As the laser beam from the laser aiming system, which has been transmitted through the wedge prism, is directed onto the position detecting element as the target, the position of a laser spot can be detected, the deflection angle at the point where the wedge prism is arranged is measured by the quantity of turning of the wedge prism, and the distance between the individual measuring points is also measured by the range finder of the wedge prism. According to this conventional technique, the position of excavation by the underground excavator is measured in terms of coordinate positions on the basis of the distances between the respective measuring points, the deflection angle and the position of the laser spot obtained as described above.\nAccording to the former conventional apparatus, the laser beam oscillator is caused to turn in the yawing direction and/or the pitching direction such that a laser beam, which is a laser light of a high convergence degree, is irradiated to a predetermined position on a laser beam receiver. The turned angle or angles are detected, and based on the turned angle or angles so detected, the position of the shield machine offset from the planned line is computed and determined by a computer. Upon measuring the position of excavation by the underground excavator, an operation for turning the laser beam oscillator to precisely direct a laser beam to the predetermined position on the laser beam receiver is therefore needed. This makes the operation complex and moreover, requires a turning mechanism for turning the laser beam oscillator, such as servomotors, thereby also making the mechanism complex. As a consequence, various problems have arisen. For example, because of the provision of such a turning mechanism, the apparatus becomes too large as one to be arranged in a tunnel, and needless to say, requires higher manufacturing cost. As the turning mechanism is a mechanical one, mechanical errors are added to optical errors, thereby making it difficult to assure high measuring accuracy and also making the apparatus weak against vibrations.\nOn the other hand, the deflection angle measuring technique which is adopted in the latter conventional technique is a method in which a deflection angle is optically measured by a laser beam. Of the problems observed on the above-mentioned method making use of a transit, the problem that measurement errors occur due to mechanical measurement errors by the turning mechanism and external force such as vibrations therefore still remains unsolved, although the problems caused by the dependence on human ability, such as the need for labor and the requirement of a long time for each measurement, have been reduced. Described specifically, according to the deflection angle measuring technique making use of this conventional technique, to always direct a laser beam from the laser aiming system, said laser beam being a laser light of a high convergence degree, onto the position detecting element as a target by a turning operation of the wedge prism, the turning mechanism is required. Like the method making use of a transit, mechanical measurement errors tend to occur due to the turning mechanism, thereby making it difficult to assure high measuring accuracy.\nFurthermore, according to the deflection angle measuring technique making use of this conventional technique, an application of external force, such as that causing tilting in the yawing direction or the pitching direction like vibrations, to the laser aiming system, the wedge prism or the target results in a measurement error by the external force, so that the measurement results of a deflection angle are affected significantly. This external force, such as vibrations, is applied by vibrations of street traffic or construction work, wind force or the like on the ground or by yawing or pitching of an underground excavator or vibrations due to excavation by the underground excavator or falling of earth in the ground. Especially in the measurement of the position of excavation by the underground excavator, the measurement results of a deflection angle affects the measurement results of the position of excavation to a considerable extent and moreover, there is a high degree of need for the accurate measurement of deflection angles because such deflection angles are often measured at gently curved places. When a measurement error takes place due to a mechanical measurement error associated with the turning mechanism or external force such as vibrations, the measurement results of the position of excavation by the underground excavator are affected considerably.\nThe present invention is intended to eliminate these problems which are observed on the conventional techniques, and has as a first object the provision of a position measuring apparatus for an underground excavator, which upon measuring a position of excavation by the underground excavator, does not need an operation for directing a light onto light-receiving means and hence does not require any mechanism for such an operation.\nFurther, the present invention also has as a second object the provision of an optical deflection angle measuring apparatus which upon optically measuring a deflection angle, does not require an operation for directing a light onto a position detecting element and prevents external force such as vibrations, even when applied, from affecting the measurement results.\nThe first object of the present invention is achieved by a position measuring apparatus for an underground excavator, said apparatus being useful in measuring a position of excavation by the underground excavator advancing in the ground while excavating a tunnel and being adapted to measure a position of a measured point, which is arranged at a front in an excavating direction and serves as an indication for the position of excavation, on a basis of its positional relationship with a measuring cardinal point arranged at a rear in the excavating direction and serving as a cardinal point for the measurement, comprising:\na cardinal-point measuring unit for setting the measuring cardinal point, said cardinal-point measuring unit having a light source capable of emitting a diffuse light forward, converging means capable of converging a diffuse light from a forward light source, and light-receiving means arranged such that the light-receiving means can receive a light converged by the converging means and can detect a direction of the forward light source on a basis of a position of the light so received;\na measured-point measuring unit for setting the measured point, said measured-point measuring unit having a light source capable of emitting a diffuse light rearward, converging means capable of converging a diffuse light from a rear light source, and light-receiving means arranged such that the light-receiving means can receive a light converged by the converging means and can detect a direction of the rear light source on a basis of a position of the light so received;\nat least one intermediate measuring unit arranged between the cardinal-point measuring unit and the measured-point measuring unit in the tunnel and having a light source capable of emitting diffuse lights forward and rearward, converging means capable of converging diffuse lights from forward and rearward light sources, respectively, and light-receiving means arranged such that the light-receiving means can receive respective lights converged by the converging means and can detect directions of the forward and rearward light sources on a basis of positions of the respective lights so received; and\ncomputing means for computing and measuring a position of the measured point relative to the measuring cardinal point on a basis of data, which are concerned with the directions of the respective light sources as determined based on detection results at the respective measuring units consisting of the cardinal-point measuring unit, the measured-point measuring unit and the intermediate measuring unit, and data on distances between the individual measuring units and their adjacent measuring units.\nOwing to the adoption of these technical means in the position measuring apparatus of the present invention for the underground excavator, the intermediate measuring unit converges by converging means diffuse lights from the front and rear light sources of the adjacent measuring units located at the front and rear of the intermediate measuring unit, the thus-converged lights are received by the corresponding light-receiving means, directions of the front and rear light sources relative to the intermediate measuring unit can be detected based on the positions of the respective lights so received, and data on the directions of the individual light sources can be obtained based on the detection results. In this case, the light sources capable of emitting diffuse lights are used as light sources, and the diffuse lights from the light sources are converged by the converging means and are directed onto the light-receiving means. Unlike the conventional art, it is therefore unnecessary to perform an operation for directing a light from each light source onto its corresponding light-receiving means. Once the data on the directions of the respective light sources are obtained as described above, the angles of individual advancing routes, which connect the adjacent measuring units with each other, relative to a starting direction line can be indirectly determined by computation without directly detecting them from their relationships with the starting direction line. In this case, even when the intermediate measuring unit and the measured-point measuring unit deviate depending on their spatial positions upon mounting or deviate depending on yawing or pitching while the underground excavator is advancing, the angles can be accurately determined under conditions with such influence excluded. If data concerning the distances between the individual measuring units and their adjacent measuring units are separately collected by a suitable method, the position of the measured point relative to the measuring cardinal point can be computed and measured from the above-described data on the respective angles as determined by computation and also from the data on the respective distances.\nThe second object of the present invention is achieved by an optical deflection angle measuring apparatus for measuring with lights a deflection angle between two line segments connecting a cardinal point to respective points set on opposite sides of the cardinal point with distances left from the cardinal point, comprising:\ndeflection-angle measuring light sources for emitting diffuse lights, said deflection-angle light sources being arranged at the respective points set on the opposite sides of the cardinal point; and\na deflection-angle measuring detector arranged at the cardinal point and provided with:\na common converging means for converging at least portions of the diffuse lights from the respective light sources,\nposition detecting elements for receiving the lights, which have been emitted from the respective light sources and have been converged by the converging means, and detecting received positions of the converged lights, and\nlight-direction changing means for changing direction of at least portions of the diffuse lights from the respective light sources, which are about to enter the converging means, such that the at least portions of the diffuse lights are allowed to be transmitted and lights, which have been emitted from the respective light sources and have then been converged by the converging means, are guided to the corresponding position detecting elements;\nwherein the respective position detecting elements are arranged at positions where the position detecting elements do not block the diffuse lights from the respective light sources, said diffuse lights being about to enter the converging means, whereby the deflection-angle measuring detector is constructed and a deflection angle between optical axes of the respective light sources can be calculated by computation on a basis of results of detections by the respective position detecting elements.\nThe optical deflection angle measuring apparatus according to the present invention, upon optically measuring a deflection angle, does not require an operation for directing a light onto each position detecting element unlike the case making use of laser beams as light sources, because as described above, the light sources capable of emitting especially diffuse lights are used as the deflection angle measuring light sources, which are arranged at the respective points set on the opposite sides of the cardinal point, to permit illumination of wide areas with spread lights and coupled with this, the spread lights from the respective light sources are converged by the converging means and received at the corresponding position detecting elements to determine the received positions of the lights and the deflection angles between the optical axes of the individual light sources can be obtained based on the respective detection results. Further, the light direction changing means are arranged to transmit at least portions of the diffuse lights from the respective light sources, said diffuse lights being about to enter the converging means, and also to change the directions of the lights form the respective light sources, said lights being about to be converged by the converging means, such that the thus-converged lights are guided to the corresponding position detecting elements, and the individual position detecting elements are arranged at positions where the, individual position detecting elements do not block the corresponding diffuse lights which are about to enter the converging means. The incidence of the diffuse lights from the corresponding light sources into the converging means is therefore not prevented by the corresponding position detecting elements. It is accordingly not required to provide each light source with its own converging means for converging a diffused light, and a common converging means can be used for the light sources. Owing to these features, even when external force such as vibrations in the pitching or yawing direction is applied to the deflection angle measuring detectors or light sources, the measurement results of the deflection angle are not affected by such external force. Moreover, the measurement results of the deflection angle are not affected by the spatial positions of the mounted detectors and light sources insofar as the setting of positions is accurately conducted upon mounting the detectors and light sources."}
{"text": "The present invention relates to an unstable laser resonator with output coupler having radially variable reflectivity.\nThe main problem with laser resonators is that they generate a high-intensity beam with minimum divergence. Present resonators for high-intensity laser beams belong to the unstable resonator category.\nA known unstable laser resonator requires the employment of a totally reflecting concave mirror at one end of the resonating cavity and a small totally reflecting convex mirror at the other end. Such a resonator is characterized at the smaller mirror by a stepped reflectivity curve which causes emission of a laser beam to be null in the centre and maximum on the periphery distributed according to a circular ring with extension equal to the difference in the dimensions of the two mirrors. The output beam is thus perturbed by the diffraction effects which occur at the side edge of the smaller mirror."}
{"text": "Cellulose fibrous structures, such as paper towels, facial tissues, and toilet tissues, are a staple of every day life. The large demand and constant usage for such consumer products has created a demand for improved versions of these products and, likewise, improvement in the methods of their manufacture. Such cellulosic fibrous structures are manufactured by depositing an aqueous slurry from a headbox onto a Fourdrinier wire or a twin wire paper machine. Either such forming wire is an endless belt through which initial deterring occurs and fiber rearrangement takes place.\nAfter the initial formation of the web, which becomes the cellulosic fibrous structure, the papermaking machine transports the web to the dry end of the papermaking machine. In the dry end of a conventional papermaking machine, a press felt compacts the web into a single region cellulosic fibrous structure prior to final drying. The final drying is usually accomplished by a heated drum, such as a Yankee drying drum.\nOne of the significant aforementioned improvements to the manufacturing process, which yields a significant improvement in the resulting consumer products, is the use of through-air drying to replace conventional press felt dewatering. In through-air drying, like press felt drying, the web begins on a forming wire, which receives an aqueous slurry of less than one percent consistency from a headbox. Typically, initial dewatering takes place on the forming wire. The forming wire is not typically exposed to web consistencies of greater than 30 percent. From the forming wire, the web is transferred to an air pervious through-air-drying belt.\nAir passes through the web and the through-air-drying belt to continue the dewatering process. The air passing the through-air-drying belt and the web is driven by vacuum transfer slots, other vacuum boxes or shoes, predryer rolls, etc., and molds the web to the topography of the through-air-drying belt, increasing the consistency of the web. Such molding creates a more three-dimensional web, but also causes pinholes, if the fibers are deflected so far in the third dimension that a breach in fiber continuity occurs.\nThe web is then transported to the final drying stage where the web is also imprinted. At the final drying stage, the through-air drying belt transfers the web to a heated drum, such as a Yankee drying drum for final drying. During this transfer, portions of the web are densifted during imprinting, to yield a multi-region structure. Many such multi-region structures have been widely accepted as preferred consumer products. An example of an early through-air-drying belt which achieved great commercial success is described in commonly assigned U.S. Pat. No. 3,301,746, issued Jan. 31, 1967 to Sanford et al.\nOver time, further improvements became necessary. A significant improvement in through-air-drying belts is the use of a resinous framework on a reinforcing structure. This arrangement allows drying belts to impart continuous patterns, or, patterns in any desired form, rather than only the discrete patterns achievable by the woven belts of the prior art. Examples of such belts and the cellulosic fibrous structures made thereby can be found in commonly assigned U.S. Pat. Nos. 4,514,345, issued Apr. 30, 1985 to Johnson et al.; 4,528,239, issued Jul. 9, 1985 to Trokhan; 4,529,480, issued Jul. 16, 1985 to Trokhan; and 4,637,859, issued Jan. 20, 1987 to Trokhan. The foregoing four patents are incorporated herein by reference for the purpose of showing preferred constructions of patterned resinous framework and reinforcing type through-air-drying belts, and the products made thereon. Such belts have been used to produce extremely commercially successful products such as Bounty paper towels and Charmin Ultra toilet tissue, both produced and sold by the instant assignee.\nAs noted above, such through-air-drying belts used a reinforcing element to stabilize the resin. The reinforcing element also controlled the deflection of the papermaking fibers resulting from vacuum applied to the backside of the belt and airflow through the belt. The early belts of this type used a fine mesh reinforcing element, typically having approximately fifty machine direction and fifty cross-machine direction yarns per inch. While such a fine mesh was acceptable from the standpoint of controlling fiber deflection into the belt, it was unable to stand the environment of a typical papermaking machine. For example, such a belt was so flexible that destructive folds and creases often occurred. The fine yarns did not provide adequate seam strength and would often bum at the high temperatures encountered in papermaking.\nYet other drawbacks were noted in the early embodiments of this type of through-air-drying belt. For example, the continuous pattern used to produce the consumer preferred product did not allow leakage through the backside of the belt. In fact, such leakage was minimized by the necessity to securely lock the resinous pattern onto the reinforcing structure. Unfortunately, when the lock-on of the resin to the reinforcing structure was maximized, the short rise time over which the differential pressure was applied to an individual region of fibers during the application of vacuum often pulled the fibers through the reinforcing element, resulting in process hygiene problems and product acceptance problems, such as pinholes.\nA new generation of patterned resinous framework and reinforcing structure through-air-drying belts addressed some of these issues. This generation utilized a dual layer reinforcing structure having vertically stacked machine direction yarns. A single cross-machine direction yarn system tied the two machine direction yarns together.\nFor paper toweling, a coarser mesh, such as thirty-five machine direction yarns and thirty cross-machine direction yarns per inch, dual layer design significantly improved the seam strength and creasing problems. The dual layer design also allowed some backside leakage to occur. Such allowance was caused by using less precure energy in joining the resin to the reinforcing structure, resulting in a compromise between the desired backside leakage and the ability to lock the resin onto the reinforcing structure.\nLater designs used an opaque backside filament in the stacked machine direction yarn dual layer design, allowing for higher precure energy and better lock-on of the resin to the reinforcing structure, while maintaining adequate backside leakage. This design effectively decoupled the tradeoff between adequate resin lock-on and adequate backside leakage in the prior art. Examples of such improvements in this type of belt are illustrated by commonly assigned U.S. patent application Ser. No. 07/872,470 filed Jun. 15, 1992, now U.S. Pat. No. 5,334,289 in the names of Trokhan et al., Issue Batch No. V73. Yet other ways to obtain a backside texture are illustrated by commonly assigned U.S. Pat. Nos. 5,098,522, issued Mar. 24, 1992 to Smurkoski et al.; 5,260,171, issued Nov. 9, 1993 to Smurkoski et al.; and 5,275,700, issued Jan. 4, 1994 to Trokhan, which patents and application are incorporated herein by reference for the purpose of showing how to obtain a backside texture on a patterned resin and reinforcing structure through-air-drying belt.\nAs such resinous framework and reinforcing structure belts were used to make tissue, such as the commercially successful Charmin Ultra noted above, new issues arose. For example, one problem in tissue making is the formation of small pinholes in the deflected areas of the web. Pinholes are strongly related to the depth that the web deflects into the belt. The depth comprises both the thickness of the resin on the reinforcing structure, and any pockets within the reinforcing structure that permits the fibers to deflect beyond the imaginary top surface plane of the reinforcing structure. Typical stacked machine direction yarn dual layer reinforcing structure designs have a variety of depths resulting from the particular weave configuration. The deeper the depth within a particular location of the weave that is registered with a deflection conduit in the resin, the greater the proclivity for a pinhole to occur in that area.\nRecent work according to the present invention has shown that the use of triple layer reinforcing structures unexpectedly reduces occurrences of pinholes. Triple layer reinforcing structures comprise two completely independent woven elements, each having its own particular machine direction and cross-machine direction mesh. The two independent woven elements are typically linked together with tie yarns.\nMore particularly, the triple layer belt preferably uses a finer mesh square weave as the upper layer, to contact the web and minimize pinholes. The lower layer or machine facing layer utilizes coarser yarns to increase rigidity and improve seam strength. The tie yarns may be machine direction or cross-machine direction yarns specifically added and which were not present in either layer.\nAlternatively, the tie yarns may be comprised of cross-machine direction or machine direction tie yarns from the upper and/or lower element of the reinforcing structure. Machine direction yarns are preferred for the tie yarns because of the increased seam strength they provide.\nHowever, this design still does not solve the problem where backside leakage may be required. Reference to the prior art teachings of backside texturing do not solve this problem either. For example, the aforementioned U.S. patent application Ser. No. 07/872,470, now U.S. Pat. No. 5,334,289, teaches the use of opaque yarns to prevent curing of resin therebelow. The resin that is not cured is washed away during the belt making process and imparts a texture to the backside of the belt. However, such a teaching further states that it is preferable the machine direction yarns be opaque because the machine direction yarns are generally disposed closer to the backside surface of the reinforcing structure than the cross-machine yarns. Such a description is not correct, however, if the machine direction yarns are used as tie yarns.\nThus, the machine direction yarn must serve either one of two mutually exclusive functions: it must either remain within the lower layer to prevent texture from going too deep into the belt, or rise out of the lower layer to tie the lower layer relative to the first layer. Compounding the problem with triple layer belts is any opaque machine direction yarns used as tie yarns will disrupt the lock-on of the resin below because such yarns intermittently are disposed on the topside of the reinforcing structure.\nAccordingly, it is an object of this invention to provide a belt which overcomes the tradeoff between high seam strength and minimal pinholing. It is further an object of this invention to provide a belt which overcomes the tradeoffs between backside leakage and low resin lock-on. The prior art has not yet provided a belt which produces consumer desired products (minimal pinholing) with a long lasting belt (high seam strength and high rigidity) and which does not lose functional components during the manufacture of the consumer product (poor resin lock-on)."}
{"text": "1. Field of the Invention\nThis invention relates to methods and apparatus for filtering signals in a subsystem including a power amplifier utilizing a bank of vibrating micromechanical apparatus.\n2. Background Art\nThe need for passive off-chip components has long been a key barrier against communication transceiver miniaturization. In particular, the majority of the high-Q bandpass filters commonly used in the RF and IF stages of heterodyning transceivers are realized using off-chip, mechanically-resonant components, such as crystal and ceramic filters and SAW devices, as illustrated in FIG. 1. Due to higher quality factor Q, such technologies greatly outperform comparable filters implemented using transistor technologies, in insertion loss, percent bandwidth, and achievable rejection. High Q is further required to implement local oscillators or synchronizing clocks in transceivers, both of which must satisfy strict phase noise specifications. Again, as illustrated in FIG. 1, off-chip elements (e.g., quartz crystals) are utilized for this purpose.\nBeing off-chip components, the above mechanical devices must interface with integrated electronics at the board level, and this constitutes an important bottleneck against the miniaturization of super-heterodyne transceivers. For this reason, recent attempts to achieve single-chip transceivers for paging and cellular communications have utilized alternative architectures that attempt to eliminate the need for off-chip high-Q components via higher levels of transistor integration. Unfortunately, without adequate front-end selectivity, such approaches have suffered somewhat in overall performance, to the point where they so far are usable only in less demanding applications.\nGiven this, and recognizing that future communication needs will most likely require higher levels of performance, single-chip transceiver solutions that retain high-Q components and that preserve super-heterodyne-like architectures are desirable.\nRecent demonstrations of vibrating beam micromechanical (xe2x80x9cxcexcmechanicalxe2x80x9d) resonator devices with frequencies in the VHF range and Q\"\"s in the tens of thousands have sparked a resurgence of research interest in communication architectures using high-Q passive devices as disclosed in the above-noted patent application entitled xe2x80x9cDevice Including A Micromechanical Resonator Having An Operating Frequency and Method of Extending Same.xe2x80x9d Much of the interest in these devices derives from their use of IC-compatible microelectromechanical systems (MEMS) fabrication technologies to greatly facilitate the on-chip integration of ultra-high-Q passive tanks together with active transistor electronics, allowing substantial size reduction.\nFIG. 2 illustrates a comparison of MEMS and SAW technologies wherein MEMS offers the same or better high-Q frequency selectivity with orders of magnitude smaller size. Indeed, reductions in size and board-level packaging complexity, as well as the desire for the high performance attainable by super-heterodyne architectures, are principal drivers for this technology.\nAlthough size reduction is certainly an advantage of this technology (commonly dubbed xe2x80x9cRF MEMSxe2x80x9d), it merely touches upon a much greater potential to influence general methods for signal processing. In particular, since they can now be integrated (perhaps on a massive scale) using MEMS technology, vibrating xcexcmechanical resonators (or xcexcmechanical links) can now be thought of as tiny circuit elements, much like resistors or transistors, in a new mechanical circuit technology. Like a single transistor, a single mechanical link does not possess adequate processing power for most applications. However, again like transistors, when combined into larger (potentially, VLSI) circuits, the true power of xcexcmechanical links can be unleashed, and signal processing functions with attributes previously inaccessible to transistor circuits may become feasible.\nThe Need for High Q in Oscillators\nFor any communications application, the stability of the oscillator signals used for frequency translation, synchronization, or sampling, is of utmost importance. Oscillator frequencies must be stable against variations in temperature against aging, and against any phenomena, such as noise or microphonics, that cause instantaneous fluctuations in phase and frequency. The single most important parameter that dictates oscillator stability is the Q of the frequency-setting tank (or of the effective tank for the case of ring oscillators). For a given application, and assuming a finite power budget, adequate long- and short-term stability of the oscillation frequency is insured only when the tank Q exceeds a certain threshold value.\nGiven the need for low power in portable units, and given that the synthesizer (containing the reference and VCO oscillators) is often a dominant contributor to total transceiver power consumption, modern transceivers could benefit greatly from technologies that yield high-Q tank components.\nThe Need for High Q in Filters\nTank Q also greatly influences the ability to implement extremely selective IF and RF filters with small percent bandwidth, small shape factor, and low insertion loss. As tank Q decreases, insertion loss increases very quickly, too much even for IF filters, and quite unacceptable for RF filters. As with oscillators, high-Q tanks are required for RF and IF filters alike, although more so for the latter, since channel selection is done predominantly at the IF in super-heterodyne receivers. In general, the more selective the filter, the higher the resonator Q required to achieve a given level of insertion loss.\nMicromechanical Circuits\nAlthough mechanical circuits, such as quartz crystal resonators and SAW filters, provide essential functions in the majority of transceiver designs, their numbers are generally suppressed due to their large size and finite cost. Unfortunately, when minimizing the use of high-Q components, designers often trade power for selectivity (i.e., Q), and hence, sacrifice transceiver performance. As a simple illustration, if the high-Q IF filter in the receive path of a communication subsystem is removed, the dynamic range requirement on the subsequent IF amplifier, IQ mixer, and A/D converter circuits, increases dramatically, forcing a corresponding increase in power consumption. Similar trade-offs exist at RF, where the larger the number or greater the complexity of high-Q components used, the smaller the power consumption in surrounding transistor circuits.\nThe Micromechanical Beam Element\nTo date, the majority of xcexcmechanical circuits most useful for communication applications in the VHF range have been realized using xcexcmechanical flexural-mode beam elements, such as shown in FIG. 2 with clampedxe2x80x94clamped boundary conditions. Although several micromachining technologies are available to realize such an element in a variety of different materials, surface micromachining has been the preferred method for xcexcmechanical communication circuits, mainly due to its flexibility in providing a variety of beam end conditions and electrode locations, and its ability to realize very complex geometries with multiple levels of suspension.\nU.S. Pat. No. 6,049,702 to Tham et al. discloses an integrated passive transceiver section wherein microelectromechanical (MEM) device fabrication techniques are used to provide low loss, high performance switches. Utilizing the MEM devices also makes possible the fabrication and use of several circuits comprising passive components, thereby enhancing the performance characteristics of the transceiver.\nU.S. Pat. No. 5,872,489 to Chang et al. discloses an integrated tunable inductance network and method. The network utilizes a plurality of MEM switches which selectively interconnect inductance devices thereby providing a selective inductance for a particular circuit.\nU.S. Pat. No. 5,963,857 to Greywall discloses an article comprising a micromachined filter. In use, the micromachined filters are assembled as part of a radio to miniaturize the size of the radio.\nU.S. Pat. Nos. 5,976,994 and 6,169,321 to Nguyen et al. disclose a batch-compatible, post-fabrication annealing method and system to trim the resonance frequency and enhance the quality factor of micromechanical structures.\nU.S. Pat. Nos. 5,455,547; 5,589,082 and 5,537,083 to Lin et al. disclose microelectromechanical signal processors. The signal processors include many individual microelectromechanical resonators which enable the processor to function as a multi-channel signal processor or a spectrum analyzer.\nU.S. Pat. No. 5,640,133 to MacDonald et al. discloses a capacitance-based, tunable, micromechanical resonator. The resonators may be selectively tuned and used in mechanical oscillators, accelerometers, electromechanical filters and other electronic devices.\nU.S. Pat. No. 5,578,976 to Yao, U.S. Pat. No. 5,619,061 to Goldsmith et al. and U.S. Pat. No. 6,016,092 to Qiu et al. disclose various micromechanical and microelectromechanical switches used in communication apparatus.\nU.S. Pat. No. 5,839,062 to Nguyen et al. disclose a MEMS-based receiver including parallel banks of microelectromechanical filters.\nU.S. Pat. Nos. 5,491,604 and 5,955,932 to Nguyen et al. disclose Q-controlled microresonators and tunable filters using the resonators.\nU.S. Pat. No. 5,783,973 to Weinberg et al. discloses a micromechanical, thermally insensitive silicon resonator and oscillator.\nThe following articles are of general interest: Nguyen et al., xe2x80x9cDesign and Performance of CMOS Micromechanical Resonator Oscillatorsxe2x80x9d, 1994 IEEE INTERNATIONAL FREQUENCY CONTROL SYMPOSIUM, pp. 127-134; Wang et al, xe2x80x9cQ-Enhancement of Microelectromechanical Filters Via Low-Velocity Spring Couplingxe2x80x9d, 1997 IEEE ULTRASONICS SYMPOSIUM, pp. 323-327; Bannon, III et al., xe2x80x9cHigh Frequency Microelectromechanical IF Filtersxe2x80x9d, 1996 IEEE ELECTRON DEVICES MEETING, San Francisco, Calif., Dec. 8-11, 1996, pp. 773-776; and Clark et al., xe2x80x9cParallel-Resonator HF Micromechanical Bandpass Filtersxe2x80x9d 1997 INTERNATIONAL CONFERENCE ON SOLID-STATE SENSORS AND ACTUATORSxe2x80x9d, pp. 1161-1164.\nAn object of the present invention is to provide a method and apparatus for filtering signals in a subsystem including a power amplifier utilizing a bank of vibrating micromechanical apparatus to allow the power amplifier to operate with less power consumption.\nIn carrying out the above object and other objects of the present invention, a method is provided for filtering signals to obtain at least one desired passband of frequencies in a subsystem including a power amplifier to allow the power amplifier to operate with less power consumption. The method includes providing a bank of micromechanical apparatus. Each of the apparatus has at least one passband of frequencies. The method also includes the step of controllably switching the bank of micromechanical apparatus to select the at least one desired passband and to substantially attenuate power outside the at least one desired passband. The method includes the step of vibrating the micromechanical apparatus corresponding to the at least one desired passband.\nThe subsystem may be an RF transmitter subsystem and the at least one desired passband may be at least one desired channel.\nEach of the apparatus may be switchable and tunable.\nEach of the devices may have a Q greater than 10 and may be greater than 5000.\nThe number of devices may number more than 10 and may even number more than 1000.\nThe apparatus may be devices, filters, circuits or signal processors.\nThe power amplifier may be a non-linear or low power amplifier.\nFurther in carrying out the above objects and other objects of the present invention, in a subsystem including a power amplifier, a signal filtering apparatus is provided which filters signals to obtain at least one desired passband of frequencies and also to allow the power amplifier to operate with less power comsumption. The apparatus includes a parallel bank of micromechanical filters. Each of the filters has a passband and is capable of handling amplified signals. The apparatus also includes a controller for selecting at least one desired filter of the bank of filters to thereby provide the at least one desired passband of frequencies.\nEach of the filters may have an input connected to a common input and each of the filters has an output connected to a common output.\nThe controller may include a decoder for controlling application of an appropriate bias voltage to the at least one desired filter.\nThe apparatus may be a high-Q, high-power filter.\nEach of the filters may be switchable and tunable.\nThe present invention uses MEMS resonators/filters to condition the output/operation of a power amplifier so as to allow it to function with much higher efficiency (and thus, much lower power). For example, MEMS tanks can be used to remove high-order harmonics caused by nonlinearity around a power amplifier, thus allowing it to operate in a nonlinear fashion.\nThe above object and other objects, features, and advantages of the present invention are readily apparent from the following detailed description of the best mode for carrying out the invention when taken in connection with the accompanying drawings."}
{"text": "This invention relates to an assembly and method suitable for optimizing allocation of resources in a multi-processor time-shared computer system.\nIn a time shared computer system, each user gets pro-rated computer resources over some interval of time. These resources are various classes of memory and disk storage and, in multi-processor systems, the available CPUs. These resources are limited to no-more than k users for k processors and not all users can be active or even contained in memory at any given time.\nCurrently, each user has a priority and is given a fixed slice of computer time in which he does his processing, and then has an inactive period of time during which other users are active. The frequency of user activation is determined by the given user\"\"s priority. There is no discipline for the optimal resource allocation in support of a set of k users, based on optimizing computer resource utilization. The problem addressed in this invention is how to optimally allocate computer resources in support of a given number of active users.\nWe have now discovered novel methodology and an assembly in realization of this methodology, which are responsive and in satisfaction of this problem.\nIn a first aspect, we disclose an assembly suitable for optimizing allocation of resources in support of k users, in a multi-processor time-shared computer system, the assembly comprising:\n(i) a resource monitor comprising:\n(a) means for inputting data comprising current job requirements;\n(b) means for operating upon the current job requirements for determining required resources;\n(c) means for storing said required resources for subsequent processing by a resource allocation;\nand\n(ii) a resource allocator connected to the resource monitor comprising:\n(a) means for inputting, from the resource monitor, said required resources; and\n(b) means for executing the required resources for generating a resource allocation for said k users.\nIn a second aspect, we disclose a method suitable for optimizing allocation of resources in support of k users, in a multi-processor time-shared computer system, the method comprising the steps of:\n1) inputting data comprising current job requirements;\n2) operating upon the current job requirements for determining required resources;\n3) storing said required resources for subsequent processing by a resource allocator;\nand\n4) executing the required resources for generating a resource allocation for said k users."}
{"text": "The present disclosure relates to polymeric materials that can be formed to produce a container, and in particular, polymeric materials that insulate. More particularly, the present disclosure relates to polymer-based formulations that can be formed to produce an insulated non-aromatic polymeric material."}
{"text": "The invention relates to a hose line having a flexible hose, having a sleeve-shaped hose connector which by way of a first hose connector part-region that is configured as a connector port projects through a hose opening that is provided at an end side of the hose into a contiguous hose end region of the hose, and which connector port has a second hose connector part-region which projects beyond the end side of the hose, which second hose connector part-region at the free end side thereof has an encircling seal retaining collar that projects on the circumference, wherein an annular flange is provided on the circumference of the second hose connector part-region, a union nut engaging behind said annular flange, and an annular seal being provided between said annular flange and the seal retaining collar.\nThe invention also relates to a filter screen for such a hose line.\nHose lines of the type mentioned at the outset are already known in the most varied embodiments. Such hose lines are thus used, for example, in order for a sanitary outlet fitting or a like water consumer point to be connected to an angle valve provided on a wall (cf. EP 2 940 364 A1). However, in sanitary water lines there is the problem that the water flow streaming through a water line can entrain remnants of chalk and dirt particles which subsequently can contaminate and impede the functioning of functional units which are provided in the onward course of the water line or at the outlet end and which, for example, regulate the maximum throughput and/or aerate the outflowing water and are intended to form the outflowing water into a homogenous, splash-free and pearly soft water jet.\nSuch functional units which are configured, for example, as jet regulators, optionally having a flow regulator provided upstream, or as jet aerators and are capable of being assembled on the water outlet of a sanitary outlet fitting are typically equipped with in inflow-side filter screen in order for dirt particles of this type to be filtered on the inflow side of said functional units. However, the space for respective functional units to be equipped with a filter screen is not available in all cases.\nTherefore, a connector screw fitment provided in the course of a water line has been developed, in which connector screw fitment an annular seal can be placed between neighboring end faces of two line portions that are to be connected by way of the connector screw fitment, said annular seal in the annular opening of the annular shape thereof having a screen which by way of a crimped shell region is held in a friction-fitting or form-fitting manner on the interior circumference of the annular seal that delimits the annular opening.\nFunctional units which are configured as backflow preventers and are capable of being assembled in the course of a water line are thus already known from DE 297 03 335 U1 and DE 10 2016 012 646.9. The previously known backflow preventers have in each case one filter screen which is capable of being releasably latch-fitted to the inflow end side of said backflow preventers. However, the filter screens herein are not capable of being used without the assigned backflow preventer."}
{"text": "The present invention relates to a method and apparatus for setting the preload between die plates such as found in molding, stamping and pressing machines. More particularly, the invention is concerned with setting the preload between two die plates that are operated by a clamping mechanism of given stroke.\nIn injection molding machines of the prior art, a pair of die plates are reciprocated relative to one another to open and close molds defined or carried by the plates. A toggle linkage is customarily employed to develop the high clamping forces holding the die plates together during the injection process. The linkage reacts the clamping forces through a thrust plate and a set of adjustable tie bars that extend between the thrust plate and the die plates. The clamping load is sometimes referred to as the preload because it is developed by the linkage as soon as the molds are closed.\nIt is desirable to set the \"die height\" or preload on the die plates when different molds or other changes are made in the molding machine. The preload is directly related to the amount by which the tie bars are strained when the toggle linkage places the dies in a lock-up condition. Thus, setting the length of the tie bars which are strained between the die plates and thurst plate sets the preload. The effective tie bar length can be set by means of a motor which rotates nuts on the ends of the tie bars and moves the thrust plate relative to the die plates. U.S. Pat. No. 2,498,264 to Goldhard discloses a set of motor-driven nuts of this type and a control mechanism for deenergizing the motor when a particular setting is obtained.\nThe strain of the tie bars associated with the preload is a small quantity compared to the relative displacement of the die plates in a molding operation; therefore, the mechanism for adjusting that strain must be accurate. In the Goldhard patent referenced above, a limit switch that can be adjusted in position for various preloads is actuated by one of the die plates as the tie bar length for a particular operation is reached. Since the adjustment in length must be set to thousandths of an inch, the limit switch actuated by the die plate, thereby measuring absolute position of the plate, makes accuracy and repeatability difficult to achieve.\nIt is accordingly a general object of the present invention to provide method and apparatus for setting the preload or die height between two clamped members in molding machines and the like."}
{"text": "1. Field of Invention\nThe invention relates to a stator coil module, a method of manufacturing the stator coil module, and an electric rotating machine. In particular, the invention relates to an electric rotating machine mounted in a vehicle, a stator coil module employed in the electric rotating machine, and a method of manufacturing the stator coil module.\n2. Description of the Related Art\nAn art for cooling a stator in an electric rotating machine (a motor, a generator, or a motor-generator) is disclosed, for example, in Japanese Patent Application Laid-Open No. 2002-186211.\nIn the electric rotating machine disclosed in the aforementioned publication, a plate is provided between segmental cores, and the plate is sealed with resin. It is difficult to perform an operation of molding the plate with resin.\nIn the electric rotating machine disclosed in the aforementioned publication, a slot serves as a cooling liquid passage, and a coil is provided in the cooling liquid passage. Thus, the number of combinations of a cooling liquid flowing through the cooling liquid passage and a material for the cores is limited."}
{"text": "The present invention relates to a method and to an apparatus for calibrating tubular articles consisting of a pulverulent metallic material which has undergone isostatic compression and a subsequent sintering operation. It is to be noted that while the invention is considered to be of special value in connection with the calibration of linings for the cylinders of internal combustion engines, its general field of use is much broader and covers all tubular articles manufactured as stated above.\nIn the prior art, a tubular article of the type above specified, including cylinder linings, are calibrated by placing the lining around a cylindrical mandrel. For obvious reasons it is then necessary that the internal diameter of the lining exceeds the diameter of the mandrel by an amount sufficient to permit convenient mounting of the lining thereon. The calibration process involves subjecting the lining to radially inwardly acting forces striving at bringing the internal surface of the lining into intimate and accurate contact with the cylindrical surface of the mandrel. As is understood, this means that the internal diameter of the lining must be reduced by an amount corresponding to the initial clearance between the mandrel and the lining. This forced deformation of the lining generates inside its material reactional forces in the form of tensions tending to counteract the deformation. The tangible result of the presence of such internal stresses in the material may occur in two different forms which may be present separately or in combination."}
{"text": "1. Technical Field\nThe present invention relates to an image processing apparatus, image processing method and recording medium that are capable of extracting foreground pixels from input image data with good accuracy when performing image processing such as image compression based on layer separation.\n2. Description of Related Art\nConventionally, image formation apparatuses such as copy machines or multi-function peripherals that form images by using electronic photographic methods, inkjet methods or the like have been widely used. Moreover, recently, multi-function peripherals are becoming even more highly functionalized, with functions such as saving documents read by a scanner as electronic data, managing saved document files, transmitting document data as electronic mail and the like in demand. Documents that are read by a scanner are saved as image data, however, the size of this image data is generally large. Therefore, when storing or transmitting image data, technology for compressing images in order to reduce the amount of data is absolutely necessary.\nOne kind of image compression technology for achieving a high compression ratio is compression technology such as MRC (Mixed Raster Content) that separates an image into layers and compresses the image data. This compression technology separates the image related to the inputted image data into two layers of image data, a foreground layer and a background layer, and then compresses the image data of each layer. It is possible to increase the compression ratio of the final compressed image. The foreground layer is an image comprising text or line arts (hereafter referred to as text or the like), and compression by a lossless compression method such as JBIG (Joint Bi-level Image experts Group), MMR (Modified Modified Read) or LZW (Lempel Ziv Welch) is suitable. The background layer is an image of everything other than text or line arts, and compression by a lossy compression method such as JPEG (Joint Photographic Experts Group) is suitable.\nJapanese patent application 2005-20227 (hereafter referred to as public document 1) discloses an image compression apparatus with the object of compressing read image data of color document images that include text or photographs at a high compression ratio without a decrease in legibility of text or the like. This image compression apparatus examines the drawn state of pixels that form a drawn object such as text or graphic from image data of a color document image read by a scanner or the like, and extracts a group of connected drawn pixels as a connected pixel group. Based on the characteristics of the pixels that make up the extracted connected pixel group, the image compression apparatus classifies an image into a foreground image that is suited for a compression method that places importance on resolution, and a background image that is suited for a compression method that places importance on color tone characteristics, and then compresses each image by the most suitable compression method."}
{"text": "1. Field of the Invention\nThe present invention relates to a system having a plurality of axes oriented in different directions such as a machine tool. More particularly, the present invention is concerned with a position control method which is used in a system of the type described and which can minimize the path error while attaining high response velocity and control stability.\n2. Related Art\nA typical conventional position control method in a system of the type described relies upon a so-called soft servo system shown in FIG. 8 conducted on each of a plurality of axes. FIG. 7 shows an example of a position control system for controlling the position of an object along a circular path. The system illustrated in FIG. 7 basically employs position controls along two axes. Namely, the system shown in FIG. 7 has a position command generating means 1 capable of outputting a shape (position) command signal (R.sub.x, R.sub.y) and control means 2X, 2Y which receive the command values R.sub.x, R.sub.y of the components of the respective axes and delilver control inputs U.sub.x, U.sub.y to the control objects 3X, 3Y having the respective driving means. In operation, the control means 2X, 2Y of the respective axes, upon receipt of the respective command values R.sub.x, R.sub.y from the position command generating means 1, operate to simultaneously control the corresponding control objects 3X, 3Y, whereby a circular path is formed.\nFIG. 8 shows, for the purpose of clarification and simplicity of explanation, the detail of the position control system for uni-axial position control. This system employs an outer control loop of a comparatively low level gain (.omega..sub.o) intended for controlling the position and an inner control loop of a comparatively high level gain (.omega..sub.c) intended for controlling velocity.\nThe velocity control loop, which has a higher level of gain (.omega..sub.c) exhibits a greater stability against any disturbance and fluctuation of parameters. In addition, since the position control loop has a comparatively low level of gain (.omega..sub.o), there is no risk of imparting excessive impact (acceleration) to the mechanical system. For these reasons, the known system shown in FIGS. 7 and 8 offers an advantage in that no specific consideration is needed in the formation of the machining or processing program when this sytem is applied to, for example, a numerical control machine tool.\nThis known position control system however, still suffers from the following problems.\nNamely, this known system does not enable the response speed and the stability of the system to be adjusted independently. Thus, the system is often required to operate under such a compromise that either one or both of the response speed and the stability do not reach the required level. Even if such a compromise could be obtained, the known system still encounters a problem in that much labor and time are required for the purpose of comparison of the instant value with the command value. Furthermore, a critical problem encountered with this known system is that the precision of the position control is undesirably limited by the maximum value of the gain (.omega..sub.o) of the position control loop.\nThis critical problem, which is serious particularly from the view point of current demand for higher accuracy or precision, will be described in more detail.\nThe characteristic (G(s)) of the position control loop in the position control system shown in FIG. 8 can be handled as a servo having a first-order lag, because the gain (.omega..sub.c) of the velocity control loop is greater than the gain (.omega..sub.o) of the position control loop. For example, the gain (.omega..sub.c) of the velocity control loop is usually 4 to 20 times greater than the gain (.omega..sub.o) of the position control loop. The characteristic (G(s)) therefore is represented as follows. ##EQU1##\nOn the other hand, the response characteristic V(t) to a stepped velocity command (V.sub.0 /S) is represented as follows. EQU V.sub.(t) =V.sub.o (1-e.sup.p) (2)\nwhere, P represents -.omega..sub.o t. PA1 =( .sub.x, .sub.y) . . . increment of position command vector PA1 = - =(e.sub.x, e.sub.y) . . . position error vector PA1 =(U.sub.x, U.sub.y) . . . control input vector PA1 =(C.sub.x, C.sub.y) . . . position vector PA1 = + PA1 = PA1 : state variable PA1 : coefficient matrix of control object PA1 : input matrix\nTherefore, the acceleration a.sub.(t) required can be obtained as follows, by conducting a first order differentiation of the formula (2) as follows. EQU a.sub.(t) =V.sub.o .multidot..omega..sub.o .multidot.e.sup.p ( 3)\nA machine system which may be a machine tool is assumed here to have a maximum cutting feed velocity V.sub.0 max. In such a system, the maximum acceleration a max is given by the following formula (4). EQU a.sub.max =V.sub.o max .multidot..omega..sub.o ( 4)\nOn the other hand, the machine system itself has an allowable maximum acceleration A max which is determined by the construction thereof. Obviously, the condition of a max .ltoreq.A max has to be met. Therefore, the maximum value .omega..sub.0 max was of the gain of the position control loop is limited by the following condition. EQU .omega..sub.o max =A.sub.max /V.sub.o max (5)\nThe shape precision during cutting of a circular work is represented here in terms of a radius reduction rate .delta. of the circle. The radius reduction rate .delta. is represented as follows. ##EQU2##\nWhere, R represents the radius (mm) of the diameter, while .DELTA.R represents the decrement (mm) of the radius. V.sub.0 represents the cutting velocity (mm/sec).\nWhen a circle of a radius R is scribed at a constant cutting velocity V.sub.0, the radius reduction rate .delta. is proportional to 1/.omega..sub.o.sup.2, i.e., varies in inverse proportion to .omega..sub.o.sup.2. On the other hand, the maximum value .omega..sub.o max of the gain of the position control loop is limited due to the reason concerning the machine system, as explained before in connection with formula (5). Thus, in a position control system which is approximated by the formula (1), the upper limit of the maximum value .omega..sub.o max of the gain is determined by the allowable maximum acceleration Amax and the maximum cutting feed velocity V.sub.o max which are determined by the machine system. Therefore, when the factors (Amax, V.sub.0 max) of the machine system are definitely determined, the maximum gain .omega..sub.o max is automatically limited, thus making it impossible to obtain a higher gain. In consequence, the precision 8 is limited by the maximum value .omega..sub.o max of the gain, thus preventing attainment of any higher degree of precision.\nThe known position control method explained before in connection with FIG. 7 involves another problem. Namely, in this system, a single position command generating means 1 is adapted for conducting a machining along a circular path by conducting bi-axial position control, i.e., a position control in the direction of a first axis and a position control in the direction of a second axis. To this end, the position command generating means outputs command values R.sub.x and R.sub.y for the respective directions or axes, so that the respective control means 2X and 2Y operate so as to control the conrol systems 3X and 3Y on the basis of the commands R.sub.x and R.sub.y. Therefore, the position control by the system shown in FIG. 7 essentially requires that two control groups each consisting of the control means 2X or 2Y and the control system 3X or 3Y, have exactly the same mechanical and electrical characteristics, otherwise the control may be effected such that one of the shafts moves ahead while the other runs aback in response to the commands, resulting in an error in the determination of the path."}
{"text": "Lean-burning engines, or engines that run on an air/fuel mixture with a stoichiometrically greater amount of air than fuel, can offer improved fuel economy relative to engines configured to run on stoichiometric air/fuel mixtures.\nHowever, lean-burning engines also may pose various disadvantages. For example, burning a lean air/fuel mixture may decrease the reduction of nitrogen oxides (collectively referred to as “NOx”) in a conventional three-way catalytic converter.\nVarious mechanisms have been developed to reduce NOx emissions in lean-burning engines. One mechanism is a NOx trap. The NOx trap is a catalytic device typically positioned downstream of a catalytic converter in an emissions system, and is configured to retain NOx when the engine is running a lean air/fuel mixture and then release and reduce the NOx when the engine runs a more rich air/fuel mixture.\nA typical NOx trap includes one or more precious metals, and an alkali or alkaline metal oxide to which nitrogen oxides adsorb as nitrates when the engine is running a lean air/fuel mixture. The engine can then be configured to periodically run a richer air/fuel mixture. The nitrates decompose under rich conditions, releasing the NOx. This reacts with the carbon monoxide, hydrogen gas and various hydrocarbons in the exhaust over the precious metal to form N2, thereby decreasing the NOx emissions and regenerating the trap.\nThe use of a NOx trap can substantially reduce NOx emissions from a lean-burning engine. However, SO2 produced by the combustion of sulfur in fuel can form sulfates, which can poison the NOx storage sites and lower the NOx storage capacity of the trap.\nThe NOx storage capacity of the trap may be recovered by operating the trap for several minutes at a high temperature (for example, around 700° C.) under rich conditions. However, this process can result in the formation and emission of hydrogen sulfide, which has an unpleasant odor. The emission of hydrogen sulfide may be suppressed by alternating between lean and rich conditions while holding the NOx trap at desulfation conditions. However, this may slow desulfation significantly.\nGerman Published Patent Application No. DE 198 49 082 A1 teaches a multistage desulfation process. In the first stage, a NOx trap is exposed to slightly rich conditions (air/fuel ratio=0.98) and a relatively low desulfation temperature for a first period of time. In the second stage, the air/fuel ratio is modulated about the initial value. As the second stage progresses, the amplitude of the modulation is increased, the temperature is increased, and the frequency and midpoint of the modulation are decreased. This method may decrease the time required for desulfation relative to fixed amplitude/frequency modulation schemes. However, this method may still cause the production of excess hydrogen sulfide, and/or take more time than necessary to complete desulfation, as it does not take into account an amount of hydrogen sulfide in a trap at any instant during the desulfation process.\nThe inventors herein have recognized that the formation and emission of hydrogen sulfide during desulfation may be more efficiently addressed by utilizing a method of desulfating a catalytic NOx storage and conversion device that includes determining an amount of sulfur stored in the catalytic NOx storage and conversion device; determining an interval for exposing the catalytic NOx storage and conversion device to a rich exhaust stream based upon the determined amount of sulfur stored, wherein the interval is longer for lower amounts of sulfur stored and shorter for higher amounts of sulfur stored; and exposing the catalytic NOx storage and conversion device to the rich exhaust stream for the determined interval."}
{"text": "Field of the Invention\nThe invention relates to a method and a system for exchanging data between a program-controlled device, in particular a microcontroller, and a logic circuit."}
{"text": "In recent years, the use of motorcycles and similar two wheeled vehicles has increased tremendously. \"Bikes\" as they are referred to in the vernacular, are used for recreation, such as trailriding; sport, such as racing and hill climbing; and last, but not least, as a basic economical mode of transportation.\nBike riding, however, whether for sport, or pleasure or necessity, subjects the rider to substantial danger of personal injury because in the event of a collision or loss of control of the bike, for any reason, the odds are that the rider will be thrown from the vehicle. Statistically it has been determined that a cyclist, who is thrown, will suffer severe, if not fatal head injuries unless some type of head protection is worn. Accordingly, most States have passed legislation requiring bike operators and passengers to wear protective head gear or suffer the penalty of fine or loss of operators license.\nAs a consequence of the noted events, protective head gear has been devised and marketed by many manufacturers and in general may be described as a high impact resistant plastic helmet shell covering the head except for the facial area and provided interiorily with some type of cushioning or suspension system whereby the helmet is firmly attached to the wearer's head but in the case of accident the cushioning or suspension system prevents violent impact from being transmitted from the impact resistant shell to the head.\nConcurrently with the problem of head protection for bike operators another safety problem has evolved and that is the problem of providing adequate rear view vision for bike riders. Obviously rear vision is just as much a necessity for a bike rider as it is for the operator of a four wheeled vehicle in order to enable maneuvering of the vehicle with safety during lane changes, turning, or taking evasive action to escape possible rear end collision from a following vehicle. To date the provisions for rear vision have been unsatisfactory and in and of themselves constitute a safety hazard in the event of an accident. In general, rear vision is provided by one or more mirrors clamped to the bike handle bars and having an elongated standard which projects upwardly and outwardly to position the reflecting surface, i.e., the mirror in such a position as to provide the most wide spread field of rear vision. Since the mirrors or mirror are mounted ahead and to the side of the operator it is obvious that the operator himself or herself, as the case may be, blocks at least some of the viewing area to the rear of his trunk. Secondly, vibrations from the road and the bike engine are transmitted at least in part through the handle bars and the standard to the mirror itself with the result that, many times, the mirror or mirrors are vibrating so badly as to be substantially useless in providing any comprehensible rear vision. Additionally, in times of inclement weather exposed mirrors become so streaked or wet that the viewing area is practically obliterated.\nFinally, should the bike operator be unfortunate enough to become involved in an accident, the mirrors, particularly the standards, become potentially lethal instruments should the operator be thrown forward so that his or her body comes in contact with them."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for controlling operating of a Switched Reluctance Motor (SRM), particularly, to a method for controlling operating of a SRM to drive the SRM according to a phase under the normal condition in case a noise is inputted from the outside.\n2. Description of the Background Art\nFIG. 1 is a schematic view showing 6/4 and 12/8 structures of a conventional Switched Reluctance Motor (SRM) in accordance with the conventional art.\nAs shown in FIG. 1, a conventional SRM is composed of a SRM having 6/4 and 12/8 structures and a circuit for operating the SRM and detecting a current flowing in the respective phases applied to a winding in the SRM. Here, the circuit will be described with reference to FIG. 2 as follows.\nFIG. 2 is a circuit diagram for detecting a current applied to the conventional SRM.\nAs shown in FIG. 2, the circuit for detecting the current applied to the SRM includes switching devices Q1 to Q6 connected to respective windings LA, LB and LC in the SRM in a row, diodes D1 to D6 connected to the respective windings LA, LB and LC in a row and resisters Rd1, Rd2 and Rd3 for sensing the current flowing in the switching devices Q2, Q4 and Q6.\nOn the other hand, there are two methods for detecting the noise inputted to a position sensor in the SRM. The one method is sensing a speed change of the SRM and the other method is sensing the current flowing in the SRM. Hereinafter, the two methods will be described.\nFirst, in case of sensing the noise inputted to the position sensor in the SRM from the outside as a speed change of the SRM, the phase applied to the winding in the SRM is shifted until the SRM becomes a normal Revolutions Per Minute (RPM). At this time, the phase is shifted in the method that the B phase is turned on in case the A phase must be turned on and the C phase is turned on in case the B phase must be turned on (A phase-B phase, B phase-C phase and C phase-A phase). Namely, under the condition that the SRM is operated normally, the RPM rotates normally by excitating a magnetic phase and in case a noise is inputted, the RPM is decreased. At this time, the phase is shifted to the next phase after judging whether there is an error by sensing the RPM. Here, the process is performed repeatedly until the RPM is operated normally.\nOn the other hand, in case the noise is sensed by detecting the current applied to the SRM, a phase current flowing in respective phases applied to the winding in the SRM is detected. Whether the phases are excitated or not is judged by comparing the detected value and the standard value. In case the magnetic phase is not excitated, the SRM is operated shifting the phases one by one.\nFIG. 3 is a view illustrating in case a sensor signal is correctly recognized or in case the sensor signal is incorrectly recognized when the current applied to the SRM is detected.\nAs shown in FIG. 3, the increased amount of current is appeared to be different according to the sensor signal outputted from the position sensor for detecting the position of the rotor in the SRM. Namely, the current values in case the sensor signal is recognized correctly and in case the sensor signal is recognized incorrectly are very different. At this time, it is judged whether the sensor signal is recognized correctly or not by comparing the current value detected in operating the SRM and the standard current value. If the sensor signal is recognized incorrectly, the SRM is operated shifting the phases one by one and the operation is repeated until the sensor signal is recognized normally.\nHowever, in the conventional method for controlling operating of the SRM, the standard value must be determined sensing the speed of the SRM or the current flowing in the SRM to sense the input of the noise. The standard RPM or standard current value must be set in case the operation speed of the SRM is changed or the voltage applied to the SRM is changed by external circumstances. Particularly, in case the speed of the SRM must be changed according to the status of load, there occurs a problem that setting of the standard value (a standard RPM or a standard current value) is difficult.\nAlso, in the conventional method for controlling operating of the SRM, there occurs a problem that the SRM is stopped if inertia of load is small in case the SRM is operated shifting the phase when the noise is inputted in the SRM.\nAlso, in the conventional method for controlling operating of the SRM, the three current sensing resisters must be connected to the circuit for detecting the current applied to the SRM to sense the current flowing in the SRM and accordingly, a plurality of analogue/digital ports must be allotted thus to increase the load in a microcomputer (not shown).\nTherefore, an object of the present invention is to provide a method for controlling operating of a Switched Reluctance Motor (SRM) capable of operating the SRM according to a phase under the normal condition in case a noise is inputted to a position sensor in the SRM from the outside.\nTo achieve these and other advantages and in accordance with the purpose of the present invention, as embodied and broadly described herein, there is provided a method for controlling operating of a SRM having a rate of poles of a stator and a rotor as 6/4 or 12/8, including the steps of generating a sensor signal at every mechanical angle to sense whether the respective phases applied in the SRM are commutated and operating the SRM by consequently turning on each applied phase in SRM at the point of detecting a rising edge of the sensor signal and turning off the phase at the point of detecting a falling edge of the next sensor signal.\nThe foregoing and other objects, features, aspects and advantages of the present invention will become more apparent from the following detailed description of the present invention when taken in conjunction with the accompanying drawings."}
{"text": "When the signal light input to an optical receiving module applied in an optical communication system has large power, a large amplitude signal is input to a receiving circuit after a light receiving element such as an Avalanche Photo Diode (APD). As a result, the duty variation due to the amplitude saturation may deteriorate the received signal characteristic. In a long distance optical communication, since many optical amplifier may be provided at some points of an optical transmission line, a wide optical power range (dynamic range) is requested which can be received by an optical receiving module. On the other hand, for example, a technology has been known which has a variable optical attenuator before an internal light receiving element of an optical receiving module and attenuates signal light when input signal light has large power.\nIn recent years, an input-power monitoring function which monitors the input power of a signal light to an optical receiving module may be requested in an optical receiving module. (Refer to Japanese Laid-open Patent Publication No. 2006-203179, for example). The input-power monitoring may be required to indicate the input power of a signal light before attenuated by a variable optical attenuator even while the signal light is being attenuated. On the other hand, for example, a technology has been known in which an optical coupler before the variable optical attenuator partially splits signal light, and the split signal light is monitored for implementing the input power monitoring.\nHowever, the technology in the past may not monitor input power with high precision. For example, a branch ratio of an optical coupler may depend on the wavelength of signal light or temperature. Thus, the technology in which signal light is partially split by an optical coupler before a variable optical attenuator for monitoring may provide input-power monitoring results which vary in accordance with the changes in branch ratio of the optical coupler even for constant input power of signal light.\nOn the other hand, a look-up table may be created for wavelengths or temperatures in advance, and split signal light and the look-up table may be used to estimate the input power. However, creating a look-up table with high precision for correcting monitoring result variations due to the wavelength or temperatures may require many experiments and/or simulations, for example, which may increase the manufacturing costs for the optical receiving module.\nThe technology which partially splits and monitors signal light with an optical coupler before a variable optical attenuator may unconditionally split a part of signal light, which may cause signal light loss. Therefore, when signal light has small input power, the minimum receiving sensitivity may deteriorate."}
{"text": "1. Field of the Invention\nThe invention relates to a method for growing a thin film using a gas cluster ion beam (GCIB).\n2. Description of related Art\nGas-cluster ion beams (GCIB's) are used for etching, cleaning, smoothing, and forming thin films. For purposes of this discussion, gas clusters are nano-sized aggregates of materials that are gaseous under conditions of standard temperature and pressure. Such gas clusters may consist of aggregates including a few to several thousand molecules, or more, that are loosely bound together. The gas clusters can be ionized by electron bombardment, which permits the gas clusters to be formed into directed beams of controllable energy. Such cluster ions each typically carry positive charges given by the product of the magnitude of the electron charge and an integer greater than or equal to one that represents the charge state of the cluster ion.\nThe larger sized cluster ions are often the most useful because of their ability to carry substantial energy per cluster ion, while yet having only modest energy per individual molecule. The ion clusters disintegrate on impact with the substrate. Each individual molecule in a particular disintegrated ion cluster carries only a small fraction of the total cluster energy. Consequently, the impact effects of large ion clusters are substantial, but are limited to a very shallow surface region. This makes gas cluster ions effective for a variety of surface modification processes, but without the tendency to produce deeper sub-surface damage that is characteristic of conventional ion beam processing.\nConventional cluster ion sources produce cluster ions having a wide size distribution scaling with the number of molecules in each cluster that may reach several thousand molecules. Clusters of atoms can be formed by the condensation of individual gas atoms (or molecules) during the adiabatic expansion of high pressure gas from a nozzle into a vacuum. A skimmer with a small aperture strips divergent streams from the core of this expanding gas flow to produce a collimated beam of clusters. Neutral clusters of various sizes are produced and held together by weak inter-atomic forces known as Van der Waals forces. This method has been used to produce beams of clusters from a variety of gases, such as helium, neon, argon, krypton, xenon, nitrogen, oxygen, carbon dioxide, sulfur hexafluoride, nitric oxide, and nitrous oxide, and mixtures of these gases.\nSeveral emerging applications for GCIB processing of substrates on an industrial scale are in the semiconductor field. Although GCIB processing of a substrate is performed in a wide variety of processes, many processes fail to provide adequate control of critical properties and/or dimensions of the surface, structure, and/or film subject to GCIB treatment."}
{"text": "Systems are known for delivering medical devices, such as stents, into a body lumen. Often, such systems include a proximal portion that remains outside the body during use and a distal portion that is disposed within the body during use. The proximal portion typically includes a handle that is held by an operator of the system (e.g., a physician) during use, and the distal portion can include an outer member surrounding an inner member with a stent positioned therebetween. Generally, the operator of the system positions the distal portion within the lumen at a desired location (e.g., so that the stent is adjacent an occlusion). The operator can then retract the outer member to allow the stent to engage the occlusion/lumen wall. Thereafter, the operator removes the distal portion of the system from the lumen."}
{"text": "Electrochemical cells, or batteries, are commonly used as electrical energy sources. A battery contains a negative electrode, typically called the anode, and a positive electrode, typically called the cathode. The anode contains an electrochemically active anode material that can be oxidized. The cathode contains an electrochemically active cathode material that can be reduced. The electrochemically active anode material is capable of reducing the electrochemically active cathode material. A separator is disposed between the anode and the cathode. The battery components are disposed in a can, or housing, that is typically made from metal.\nWhen a battery is used as an electrical energy source in an electronic device, electrical contact is made to the anode and the cathode, allowing electrons to flow through the device and permitting the respective oxidation and reduction reactions to occur to provide electrical power to the electronic device. An electrolyte is in contact with the anode, the cathode, and the separator. The electrolyte contains ions that flow through the separator between the anode and cathode to maintain charge balance throughout the battery during discharge.\nThere is a growing need to make batteries that are better suited to power contemporary electronic devices such as toys; remote controls; audio devices; flashlights; digital cameras and peripheral photography equipment; electronic games; toothbrushes; radios; and clocks. To meet this need, batteries may include higher loading of electrochemically active anode and/or cathode materials to provide increased capacity and service life. Batteries, however, also come in common sizes, such as the AA, AAA, AAAA, C, and D battery sizes, that have fixed external dimensions and constrained internal volumes. The ability to increase electrochemically active material loading alone to achieve better performing batteries is thus limited.\nThe electrochemically active cathode material of the battery is another design feature that may be adjusted in order to provide increased performance. For example, electrochemically active material that has higher volumetric and gravimetric capacity may result in a better performing battery. Similarly, electrochemically active material that has a higher oxidation state may also result in a better performing battery. The electrochemically active material that is selected, however, must provide an acceptable closed circuit voltage, or running voltage, range for the devices that the battery may power. The device may be damaged if the OCV or running voltages of the battery are too high. Conversely, the device may not function at all if the running voltage of the battery is too low. In addition, electrochemically active material, such as high oxidation state transition metal oxide, may be highly reactive. The highly reactive nature of such electrochemically active material may lead to gas evolution when the electrochemically active material is incorporated within a battery and is brought into contact with the electrolyte. Any gas that is evolved may lead to structural issues within the battery, such as continuity within the cathode, and/or leakage of electrolyte from the battery. The high oxidation state transition metal oxide may detrimentally react with other battery components, such as carbon additives, e.g., graphite; other additives, e.g., surfactant(s); and/or the separator. The high oxidation state transition metal oxide may also have a tendency to consume electrolyte, which may lead to other structural issues within the battery, such as cathode swelling, and unfavorable water balance within the battery. Also, a battery including high oxidation state transition metal oxide as an electrochemically active cathode material may, for example, exhibit instability and an elevated rate of self-discharge when the battery is stored for a period of time. In addition, the ratios of both water and potassium hydroxide content to the electrochemically active material content of the battery needs to be appropriately balanced to deliver proper utilization of the electrochemically active materials. Furthermore, the selection of suitable ratios of both water and potassium hydroxide content to electrochemically active material content may provide increased battery performance across multiple discharge rate regimes.\nThere exists a need to provide a battery including an electrochemically active cathode material that address the needs discussed above. The battery including beta-delithiated layered nickel oxide electrochemically active cathode material of the present invention addresses, inter alia, these needs."}
{"text": "1. Field of the Invention\nThe present invention relates to a refrigerator, and particularly, to a refrigerator including a hot line and a plurality of capillary tubes between a condenser and an evaporator.\n2. Description of Related Art\nIn general, a refrigerator is an apparatus that keeps a storage compartment, such as a refrigerating compartment or a freezing compartment, at a low temperature by using a cooling cycle of a coolant, which consists of a compressor, a condenser, an expander, and an evaporator.\nA refrigerator may include a main body that includes a storage compartment, such as a freezing compartment and a refrigerating compartment, and a door connected to the main body to open and close the storage compartment.\nA difference in temperature occurs between the inside and outside of the refrigerator. In case the outside is left in a high-temperature, humid environment, water drops may be created on a portion of the main body, which contacts the door. This happens when room-temperature internal air transfers heat to low-temperature cold air so that the temperature of the low-temperature cold air goes up, thereby reaching the dew point. Water drops are mainly created on a region where the door contacts the main body due to a difference in temperature between the inside and outside of the refrigerator when the door is open.\nIn case a hot line is provided at the main body to be connected with the condenser, the refrigerator may remove water drops by making a coolant pass through the hot line.\nIn conventional refrigerators, a coolant which has passed through a hot line flows through one evaporator via only one expander. So, it is difficult to secure the optimum coolant circulation amount depending on load."}
{"text": "The use of metal bumps to bond external metal beam leads to contact patterns on semiconductor chips is well known in the art. Beam lead bonding, flip-chip bonding, and tape automated bonding (TAB), for example, comprise examples of semiconductor chip mounting processes that utilize metal bumps formed over selected regions of the wiring pattern on a chip. In thermal compression bonding, a combination of heat and physical pressure is used to bond metal bumps on the chip to external metal connections. Those external connections can be provided, for example, on an external chip or a flexible tape, and positioned so as to be in facing relationship to the metal bumps.\nReferring to FIG. 1, a prior art bump/chip structure is shown which has been found to present potential corrosion problems leading to chip malfunction. A semiconductor substrate 10 has formed thereon a conductive contact land 12 to which a metal bump 14 is affixed.\nA glassy, passivating layer 16 overlays the surface of substrate 10. During the processing of substrate 10, an opening is etched in layer 16 by placing a mask on its upper surface, followed by an etch. The etched opening has generally been made somewhat larger in cross section than bump 14 in order to assure that the bump, when it is deposited, is seated squarely on contact land 12.\nThe composition of bump 14 includes a bonding layer 20 of chromium on which an aluminum pedestal 22 was deposited. Next, an additional layer of chromium 24 is laid down upon the uppermost surface of aluminum pedestal 22, followed by layers of copper 26 and gold 28.\nAs a consequence of this structure exposed collar areas 18 were created where corrosion could attack the uppermost surface of contact land 12. While attempts have been made to passivate collar region 18 through the application of polymeric overcoats, the results have not been totally satisfactory.\nOther prior art showing methods for forming bumps on conductive lands on semiconductors can be found in U.S. Pat. Nos. 4,042,954, 4,427,715, and 3,874,072. In U.S. Pat. No. 4,042,954 to Harris there is shown a method for forming metal bumps which utilizes a multilayer transition structure of Cr, AlCr, Cr and Au. This transition structure is used to connect a nickel-under-copper bump to an aluminum metal pattern on a semiconductor chip. The aluminum metal pattern contacts selected regions on the chip through an SiO.sub.2 passivating layer.\nIn U.S. Pat. No. 4,427,715 to Harris, another method is shown for forming metal bumps wherein a bump is centered over a pad so as to cover a window in an intermediate passivating layer. The positioning and size of the bump is selected relative to the pads such that during thermal compression bonding, the periphery of the bump does not extend over the periphery of the pad. This arrangement is purported to prevent failures caused by cracking in the passivating layer.\nIn U.S. Pat. No. 3,874,072 to Rose et al., a method is shown for forming mushroom-shaped metal bumps which incorporates multiple layers of varying metals. Briefly, a nickel mushroom cap is bonded to an aluminum layer through intervening nickel, chromium layers respectively. The aluminum layer is deposited on an aluminum pad through a window in a thin glass passivating layer. Thin layers of gold tin and gold are formed sequentially over the nickel cap.\nThe above patents suffer from a number of disadvantages. Some do not anticipate the corrosion problems which may occur due to either a lack of a good bond between an aluminum bump and an underlying insulating structure or from exposed contact land metallurgy. Others do not anticipate that the misalignment of a bump can cause cracking of an underlying passivating layer.\nIn summary, there is a need in the art for a metal bump process, useful in thermal compression bonding which is both inexpensive to manufacture, inhibits corrosion due to exposed metallurgy, and maintains optimum sub-bump structures.\nAccordingly, it is an object of this invention to provide a new and improved metal bump for use in thermal compression bonding processes on semiconductor chips.\nIt is another object of this invention to provide an improved method for manufacturing interconnecting metal bump structures wherein passivation layer cracking is avoided.\nIt is still another object of this invention to provide a metal bump structure that does not leave regions which expose underlying contact metallurgy which is corrosion-sensitive.\nIt is yet another object of this invention is to provide a metal bump/semiconductor structure wherein a glassy passivating layer may be employed."}
{"text": "The present invention is directed to a vehicle rearview mirror system and, in particular, to a rearview mirror system incorporating one or more electro-optic reflective elements which automatically dim in response to glare-producing light.\nElectro-optic rearview mirror systems include an electro-optic reflective element and an electronic circuit, which controls a reflectance level of the electro-optic reflective element. In order to establish a reflectance level of the electro-optic reflective element, the electronic circuit includes one or more photo-sensors to sense light conditions in a particular light environment. In order to sense the desired light environment, it is common to utilize certain mounting techniques for the light sensors in order to position the light sensors to sense the appropriate light environment. The use of special mounting arrangements for the light sensors increases the cost and difficulty of assembling the electro-optic rearview mirror system."}
{"text": "This invention relates generally to a single point fuel injection system and, more particularly, to the mechanical, electromechanical and electronic portions of a fuel management system for delivering a charge of fuel to a specified opening intake valve of the engine from a single point in a throttle body.\nThe majority of automobiles being built today have fuel systems which are either controlled by means of a carburetor or a multipoint fuel injection system. While the multipoint fuel injection system has been found to be an improvement over the carburetor, it too has problems which require solution. The system being described herein is calculated to combine the advantages of both systems and either solve or ameliorate the inherent problems of the two systems.\nIn the case of a carburetor, while it has an advantage of low cost and low operating fuel pressure, there are many undesirable characteristics inherent to the use of a carburetor. For example, the operation of a carburetor provides a continuous flow of fuel, the quantity of fuel being determined by the position of the throttle. It has been found that the fuel is not properly atomized and entrained in the air flow through the throat of the carburetor. Without proper atomization, the fuel distribution to the various cylinders is uneven thereby causing a rich or lean mixture from one cylinder to another. This situation increases the objectionable emissions from the particular cylinder which is too rich or too lean relative to stoichiometric. Also, relative to a fuel injection system, the carbureted system is inherently inaccurate in its fuel control whereby all of the cylinders may be operating at a point different from optimum.\nFurther, carbureted systems are typically operated in an open loop mode of operation. With this type of operation, the output of the engine exhaust system is not sensed to determine the quality of combustion which is occuring in the engine. Under these circumstances, the optimum air/fuel ratio is not achieved and higher emission levels are again experienced.\nThe shortcomings of a carbureted system have been somewhat eliminated by certain multipoint fuel injection systems on the market. With a multipoint fuel injection system, the fuel management is provided with a rather precise control of the fuel being fed to the engine which results in improved driveability without unwanted surges, lower emission levels, convenient changes of the calibration of the system, and the system may be operated in a closed loop mode of operation.\nHowever, multipoint fuel injection systems do have certain undesirable characteristics which, if overcome, would increase the use of injected fuel management systems. For example, a typical multipoint fuel injection system involves a higher cost in the initial installation due to the sophisticated injectors being utilized and the relatively higher cost of the control electronics. Also, due to the requirement of a precise fuel pulse being fed to each cylinder, the fuel distribution between cylinders may vary due to the fact that the injectors are not perfectly matched, one to the other. As is the case with a carburetor, unless the fuel is highly atomized and rapidly carried into the appropriate cylinder immediately upon injection of fuel into the air stream, wall wetting is experienced. In the situation where the wetting of the walls with fuel is occurring, fuel is unevenly distributed to the cylinders and results in an uneven air/fuel ratio from cylinder to cylinder. Also, with wall wetting, the fuel charge being fed to the same cylinder from one cycle to the next may vary depending on the amount of fuel on the walls of the manifold. Upon injection of a fuel pulse which wets the walls of the manifold, the cylinder will receive a leaner air/fuel ratio charge than required. Subsequently, the fuel on the walls of the manifold will be entrained into the air stream to create a rich air/fuel mixture, which air/fuel mixture is not directly controlled by the duration of the fuel injection pulse. This can result in power surges which deteriorates the driveability of the automobile.\nWith a multipoint system, there are problems involved in the hot starting of the automobile and hot fuel handling due to the fact that the injectors are positioned very close to the high heat areas of the engine, as are the fuel lines feeding the injectors. This creates vaporization of the fuel resulting in a low quantity of fuel being injected per pulse to create a lean air/fuel ratio. Further, the multipoint fuel injection system requires a high pressure fuel system with the inherent sealing problems and the cost of a high pressure pump.\nWith a multipoint fuel injection system, it is seen that an injector is provided for each cylinder of the engine thereby requiring a wholly self-contained injector at each cylinder. Further, the system requires a pressure regulator which is separate from the injectors and a plurality of fuel atomizers, one for each injector being utilized in the system. It is obviously desirable to integrate all of the various parts associated with a multipoint fuel injection system into a single unit having a single housing. This reduces the cost of the system and also increases reliability."}
{"text": "1. Field of the Invention\nThe present invention relates to an annealing furnace for processing the surface of an electronic device in the manufacturing process of the electronic device and, in particular, to a manufacturing apparatus and an annealing method of forming a thin insulating film by use of a heating lamp and to a manufacturing method of an electronic device.\n2. Description of the Related Art\nLarge scale integrated (LSI) circuits have been becoming larger to improve their performance, and at the same time, finer elements have been advancing more rapidly than ever. When a metal-insulator-semiconductor field effect transistor (MISFET) is used, the electrical thickness of a gate insulating film is required to become thinner in correspondence to a decrease in a gate length to make the MISFET finer while controlling the threshold voltages thereof. Accordingly, a formation technique and optimization of thin gate insulating films have become extremely important. Especially, a technique of forming thin gate insulating films having an effective oxide thickness (EOT) of 2 nm or less is becoming increasingly important. Herein, EOT represents the thickness of the film converted into that of a silicon oxide (SiO2) film.\nFor example, a rapid thermal oxidation (RTO) or the like has been utilized as the method for forming ultra thin SiO2 or silicon oxy-nitride (SiOxNy, referred to as SiON hereinafter) films having an EOT of 2 nm or less on a semiconductor substrate such as silicon (Si), for example, (see U.S. Pat. No. 5,966,594 specification and Japanese Patent Application Laid-Open No. 6-349821). Moreover, for a method for forming an ultra thin insulating film with a thickness of 2 nm or less, an attempt has been made to use oxygen radical (O*) or oxygen ion at a low temperature of 500° C. or less.\nIf typical halogen lamps are used in the RTO, the temperature rises and falls slowly. Thus, it is difficult to form ultra thin insulating films with a thickness of 2 nm or less with great repeatability at a high oxidation temperature of, for example, approximately 1000° C. In addition, for SiON films which are formed by use of O* or oxygen ion at a low temperature, only insulating films which have high leakage current density of, for example, 100 A/cm2 under the condition of an EOT of 1 nm and an electric field of 4.5 MV/cm are obtained.\nIn a typical thermal oxidation process of Si substrates, first, a native oxide film on the surface of a Si substrate is removed by wet processing with dilute hydrofluoric acid (HF), ammonium fluoride (NH4F) or the like. However, a 0.5 nm to 1 nm thick native oxide film is formed in an atmosphere in advance of thermal oxidation. Although a Si substrate is installed in a lamp annealing apparatus in a shortest feasible time to prevent the formation of a native oxide film after wet processing, the 0.5 nm to 1 nm thick oxide film is formed by oxygen (O2), mixture gas of O2 and hydrogen (H2) or oxidation gas such as water vapor (H2O) in the atmosphere before the ambient temperature reaches the oxidation temperature.\nThe native oxide film, which is formed on the surface of the Si substrate, has inferior electrical characteristics. However, if the native oxide film is heated at, for example, 1050° C. or more, the electrical characteristics thereof are improved so that the characteristics are equivalent to those of a thermal oxide film. Thus, the thermal oxide film having good electrical characteristics can be formed in the RTO, but it is extremely difficult to control the thickness of the film to be 2 nm or less since the native oxide film is formed before the thermal oxidation process."}
{"text": "Patient monitoring devices (herein also called patient monitors), as often used in hospitals located next to the patient bed, are generally known in the art. They are used as a kind of measurement devices for measuring and displaying of physiological signs and for providing alarming and scoring based on the captured physiological signs. In case a caregiver decides to visit a patient, physiological signs, alarms and scores are available for review.\nConventional bedside patient monitors are able to measure and show the patient's physiological signs in real-time. They can also be connected to other devices, such as a PC, enabling caregivers to access and view all information available in the hospital database on a single screen.\nU.S. Pat. No. 8,145,590 B2 discloses an expert system for patient medical information analysis. A plurality of chronic sensors is used to facilitate diagnosis and medical decision making for an individual patient. An expert system evaluates the sensor data, combines the sensor data with stored probability data and provides an output signal for notification or medical intervention.\nWO 2014/091329 A1 discloses a monitoring apparatus for monitoring vital functions of a patient, comprising a connection device connectable to a patient medication storage unit and for receiving information about at least one drug administered to the patient, an interface connectable to a drug database for receiving information about at least one effect of the at least one administered drug and/or a combination of administered drugs, a measurement device associated to the patient for measuring at least one vital parameter of the patient, and a control unit for enabling a measurement of at least one vital parameter and/or for enabling an evaluation of a measurement of at least one vital parameter on the basis of the information about the at least one effect.\nV. Koutkias et al., “A personalized framework for medication treatment management in chronic care”, IEEE Transactions on Information Technology in Biomedicine, Special section on affective and pervasive computing for healthcare archive, Volume 14 Issue 2, March 2010, pages 464-472 discloses a framework for monitoring the patient's condition and safety with respect to the medication treatment administered. For this purpose, considering a body area network (BAN) with advanced sensors and a mobile base unit as the central communication hub from the one side, and the clinical environment from the other side, an architecture is described, offering monitoring patterns definition for the detection of possible adverse drug events and the assessment of medication response, supported by mechanisms enabling bidirectional communication between the BAN and the clinical site."}
{"text": "There are presently many different types of packages on the market which are useful for holding a plurality of articles in a fixed orientation with respect to each other. For example, in the cigarette industry, the most common type of package which is employed is a folding carton which completely surrounds all four side walls and the two end walls of the package. Such a folding carton thereby requires a great deal of paperboard material. For example, in packaging 100 millimeter cigarettes, the blank size required to produce such a folding carton is approximately 175 square inches. However, the increasing costs of paperboard material have created a demand for a package which would accomplish the same purpose as the above-described folding carton but which would utilize substantially less paperboard material.\nIn addition, there are other inefficiences created by the folding-type carton presently employed in the cigarette industry. That is, after the folding carton is originally packed and sealed, during its distribution cycle, it must be reopened for further processing, such as tax stamping or price marking, and then resealed. Accordingly, it would be highly desirable to provide a package which also eliminated the inefficient steps of reopening and reclosing the package for tax stamping or the like, in addition to substantially reducing the paperboard material utilized.\nThere are also other industries which utilize inefficient and costly packages. For example, the multipacking of cans and plastic and glass containers, particularly those which have an irregular shape, is frequently accomplished through the use of a paperboard overwrap which completely surrounds the four side walls and two end walls of the package and which serves to hold the articles in a fixed orientation with respect to each other. It has been realized that there would be a great economic saving if such irregular-shaped articles could be shrink-wrapped with plastic instead of overwrapped with the more costly paperboard material. However, such irregular-shaped articles are difficult to shrink-wrap because they tend to change position and move out of alignment upon being subjected to the pressure of shrink-wrapping. Therefore, as such irregular-shaped articles are difficult to maintain in a fixed orientation with respect to each other, it has prevented such industries from taking advantage of the economic savings of shrink-wrapping. It would therefore be highly desirable to provide an inexpensive manner for maintaining such irregular-shaped articles in a fixed orientation with respect to each other so that they may be subjected to the pressure of shrink-wrapping and thereby avoid the higher costs of overwrapping the articles with paperboard material.\nThis problem also extends to the case-packing of irregular-shaped articles which must be maintained in a fixed orientation with respect to each other in order to be packed into a rectangular-shaped case. Typically, an outer paperboard sleeve must be utilized to hold the articles in the desired orientation for case-packing, even though the paperboard sleeve is normaly discarded when the case is unpacked for shelf stocking in stores. Again, it would also be highly desirable to provide a simple and less costly manner of maintaining the orientation of such irregular-shaped articles for case-packing which would not require the use of he more costly paperboard sleeves.\nBroadly, it is an object of the present invention to provide an improved package which overcomes one or more of the aforesaid problems. Specifically, it is within the contemplation of the present invention to provide an improved package which is constructed from a minimum of paperboard material and still accomplishes the function of maintaining a plurality of articles, including irregular-shaped articles, in a fixed orientation with respect to each other so that such articles may be overwrapped, shrink-wrapped, case-packed, or the like.\nIt is a further object of the present invention to provide an improved package which substantially reduces material and labor costs and also increases the rate of packaging.\nIt is a still further object of the present invention to provide an improved package which maintains a plurality of articles in a fixed orientation with respect to each other, but leaves the end walls of the package exposed for further processing, such as tax stamping or price marking, and thereby eliminates the need for the inefficient steps of reopening and resealing packages.\nIt is also an object of the present invention to provide an inexpensive way of maintaining a plurality of irregular-shaped articles in a fixed orientation with respect to each other so that they may be shrink-wrapped or case-packed and thereby avoid the use of more costly paperboard overwraps or sleeves for such irregular-shaped articles."}
{"text": "1. Field of the Invention\nThe invention pertains to the field of solvent systems. More particularly, the invention pertains to a solvent system including at least one terpene.\n2. Description of Related Art\nThe U.S. federal government has assigned to the Environmental Protection Agency (EPA) the responsibility of regulating man-made volatile emissions, which may pollute the air and atmosphere. The EPA has listed a group of VOCs as “Exempt VOCs”, which is published in the Federal Register under 40 CFR 51.100(s), which is hereby incorporated by reference herein, and addendums. All other VOCs are restricted for use and weight-limited for emission discharge into the atmosphere and environment. Formulators and compounders of products using VOCs are challenged to use and obtain the required VOC combinations for use in their products, to comply with toxic emissions, and to provide for safety in manufacturing, shipping, storage, and ultimate use. VOCs are further sub-classified based upon their vapor pressures, boiling points, and flash points. The flash point of the VOC becomes a critical factor for emission discharge and safety, because it is critical because it is the point at which the liquid becomes a volatile vapor, mixes with oxygen, and thereby acquires its most combustible or flammable state. Flash point, as used herein, refers to the lowest temperature at which a volatile liquid can vaporize to form an ignitable mixture in air. In the United States, the Environmental Protection Agency (EPA) and the Department of Transportation (DOT) have classified such compounds based on their volatilities or “flash points”.\nThe current EPA and DOT Volatile Organic Compound (“VOC”) classifications are as follows:\nClass I liquids (flammable)flash point at or below 100° F.Class II liquids (combustible)flash point from 100° F. to below 140° F.Class III liquids (combustible)flash point above 140° F. to below 200° F.\nInternational standards are generally stricter, with 43° C. (109.4° F.) generally being the separation point between flammable and combustible liquids.\nThere are currently six commonly used exempt VOCs on the U.S. federal list of acceptable VOCs. Of the six, four have flash points below 10° C./50° F. and they are therefore classified as “flammable and hazardous”, thus rendering their respective use to safe and unpopulated areas. P-Chlorobenzotrifluoride (PCBTF) has a “combustible” flash point rating of 43° C./109° F., and propylene carbonate lists as “combustible” with a flash point of 107.8° C./226° F.\nOrganic solvents, such as acetone, xylene, ketones, esters, ethers, aliphatic hydrocarbons, and aromatic hydrocarbons such as benzene, are widely used as additives for industrial and commercial purposes. Due to the high volatility of these compounds, their uses are regulated by many countries' governmental agencies. Obviously, the more flammable a solvent, the more restrictions exist on its use. Further, manufacturers that utilize solvents must handle the more flammable liquids more carefully and have to address issues involving atmospheric volatility and worker health concerns due to excessive exposure to these chemicals. If these solvents can be modified so that their flash points can be increased, this would result in significantly more uses for these compounds as well as increasing their shelf lives. It is desirable therefore to increase the flash point of a variety of solvents without substantially increasing the cost of the solvent while maintaining the solvent and its utility and effectiveness for its purpose.\n“Green” solvents are regarded as such because of their sources of origin (non-petroleum based) and the fact that they biodegrade readily after use without environmental damage. Green solvents include, but are not limited to, certain alcohols such as methanol, ethanol, benzyl alcohol, certain acetates, certain esters, and turpentine. The problem with these solvents, however, is that due to their high volatility, they are considered hazardous air pollutants which violate Federal and State emissions regulations. Ethyl lactate, one example of a green solvent, has a flash point of about 115° F., which despite not being considered a Class I VOC, is still combustible, which limits its potential uses. Green esters, such as N-butyl propionate, ethyl lactate, and methyl soyate, are also considered combustible, with flash points ranging from 115° F. up to 250° F. However, they are slow to evaporate and exhibit good solvating characteristics when used in blended compositions. Nevertheless, this group of solvents is considered environmentally clean and could be utilized more extensively if their volatilities could be reduced.\nThe growing environmental concerns over the use of petroleum lubricants for industrial and transportation needs has expanded the uses of oleochemicals, vegetable oils, and recyclables for use as alternative lubricants. Environmental concerns and efforts toward resource renewability, biodegradability, and utility have led to the introduction and use of vegetable oils and synthetic esters. Of the major candidates considered, coconut oil, hydro-processed mineral oil, rapeseed (canola) oil, soy methyl esters, palm oil, and recyclable vegetable oils have been found to offer the best overall performance.\nThe art of selecting and using various oils for lubricants is routinely practiced. However, environmental regulations now restrict the use of volatile organic compounds (VOCs), which affects the lubricating oil selection and/or the method of application. U.S. federal and state regulations governing the composition and emissions content of lubricants conforms and limits choices.\nTerpenes are a large and varied class of hydrocarbons with the molecular formula (C5H8)n and are produced by a wide variety of plants and trees, such as conifers and pines. Terpenes are derived biosynthetically from units of isoprene, which has the molecular formula C5H8. The term “terpene” is sometimes used broadly to also include the terpenoids which are terpenes that have been chemically changed or modified, such as through oxidation. A terpene, as used herein, may be any hydrocarbon, natural or synthetic, formed from isoprene units. Terpenes as well as terpenoids, are the primary ingredients of the essential oils of many types of trees, plants, and flowers, including citrus fruits. Terpenes are the major components of rosin (resin) as well as turpentine produced from gum rosin (resin).\nThe terpenes are generally classified sequentially by the number of isoprene units they contain as hemiterpenes (one isoprene unit), monoterpenes (two isoprene units), sesquiterpenes (three isoprene units), diterpenes (four isoprene units), sesterterpenes (five isoprene units), triterpenes (six isoprene units), and tetraterpenes (eight isoprene units).\nWhen a terpene is modified chemically, such as by oxidation or rearrangement of the carbon skeleton, the resulting compound is generally referred to as a terpenoid. Terpenes and terpenoids are typically derived from plants, trees, flowers, and other vegetation. They come in the form of liquids, solids, waxes, oils, and alcohols. Terpenes and terpenoids may be formed as acyclic, monocyclic, or polycyclic structures. Sometimes alternatively referred to as “isoprenoids”, terpenoids are derived from five-carbon isoprene units and can be classified as “modified terpenes”, where methyl groups have been moved or removed, or oxygen atoms added. Some researchers use the term terpene to include all terpenoids.\nTerpenes and terpenoids in various forms have been used for centuries in fragrances due to their compatibility with other compounds and their minimal negative environmental impact. Terpineol, a terpene alcohol, has the chemical formula: C10H18O and is found in three isomeric forms, alpha, beta, and gamma, with beta-terpineol being non-naturally occurring. Terpenes and terpenoids have been used for other purposes, such as disinfectants, cleaning compounds, soaps, cosmetics, and colognes. They are also known to add, enhance, or mask the odor of products which might be offensive to humans or animals.\nThe term “terpene” is derived from the word “turpentine”. Turpentine is a volatile fluid which is distilled and refined for further commercial uses. The most common terpenes obtained from rosin (resin) distillation are the bicyclic terpenes alpha-pinene (α-pinene), beta-pinene (β-pinene), delta-3 carene (δ-3 carene), and sabinene, the monocyclic terpenes limonene and terpinolene, and smaller amounts of tricyclic sesquiterpenes longifolene, caryophyllene, and delta-cadinene. Rubber, which is a polyterpene, is one of the most widely known terpenes.\nAs with other plant essential oils, terpenes are major constituents of the essential oils of citrus fruits. However, they are removed before the essential oil is used for flavoring beverages and foods because they tend to produce undesirable tastes when permitted to oxidize and polymerize.\nTerpene hydrocarbons are liquid distillates separated from rosin pitch or sap from conifers, pine trees, citrus, and varied vegetation. The volatile terpene hydrocarbons, Chemical Abstracts Service Number (CAS#) 8006-64-2, are non-oxygenated with the molecular formula C10H16. Common names include terpenes, diterpenes, pure gum turpentine (PGT), oil of turpentine (OT), and limonene. Commercial areas of use include in paint thinners, paint strippers, cleaners, disinfectants, and pharmaceuticals.\nPure gum turpentine is a complex formulation made of α-pinene (CAS#80-56-8) 40%-55% weight, β-pinene (CAS#127-91-3) 25%-35% weight, and the balance mixed trace terpenes. Pure gum turpentine is 99.5%-100% volatile, evaporating slightly less (0.9) than the reference standard (butyl acetate (1.0)) and has a flash point of 35° C./95° F., a boiling point of 155-180° C., and a Kauri-butanol (Kb) value greater than 50 and is generally soluble in organic solvents.\nU.S. Pat. No. 7,273,839, issued Sep. 25, 2007 to Koetzle and hereby incorporated by reference herein, discloses the use of terpene alcohols with organic solvents and blends of solvents to increase the flash points of these solvents. Koetzle also discloses a method to decrease the flammability of normally flammable solvents using terpene alcohols, by blending the terpene alcohol into the flammable solvent. Koetzle discloses increasing the flash points of acetone, methanol, ethyl acetate, ethanol, and xylene by 50 to 60° C., by addition of 12-14% terpineol. Koetzle also discloses solvents blended with other organic solvents to produce performance solvents, such as paint strippers with flash points greater than 140° F.\nTurpentine, which is also known as spirit turpentine, oil of turpentine, and wood turpentine, is obtained by the distillation of resin from trees, usually pine trees. Turpentine prepared in this manner includes mostly alpha-pinene and beta-pinene, two terpene isomers. The exact composition of turpentine may vary from batch to batch depending on the distillation conditions and the resin source. As a solvent, turpentine has been used in varnishes and for thinning oil-based paints.\nTerpene alcohols generally have the structure of terpenes except that they include at least one hydroxyl group. A terpene alcohol, as used herein, may be any compound, natural or synthetic, formed from isoprene units and having at least one hydroxyl group. Terpene alcohols are also derived from plants, trees, flowers, and other vegetation which allows their classification as “green compounds”. Terpene alcohols are also divided into groups determined by the number of carbon atoms and repeating isoprene units. Terpene alcohols may be formed as acyclic, monocyclic, or polycyclic structures. Terpineol, a terpene alcohol, has the chemical formula: C10H18O and is found in three isomeric forms, alpha (α), beta (β), and gamma (γ), with β-terpineol being non-naturally occurring. Terpene alcohols have been used for many purposes, including, but not limited to, disinfectants, cleaning compounds, soaps, cosmetics, and colognes. They are also known to add to, enhance, or mask the odor of products which might be otherwise offensive to humans or animals.\nLatex paints are currently being used almost exclusively today, because they are water based and nonflammable and they do not emit any harmful vapors into the atmosphere. Latex paints, however, do have drawbacks in comparison to oil-based paints, which are currently banned in many places for their flammability and vapor emissions. Oil-based coatings, such as vinyl coatings and oil-based paints, are more durable than latex paints. Lines on streets and roads made with latex paints must be repainted much more frequently than lines made with oil-based paints. Oil-based coatings are also easier to spread, easier to touch up, and can be made more resistant to fungus and molds than latex paints. There is a need in the art for an oil-based coating with reduced flammability, which meets stricter government regulations."}
{"text": "1. Related Patents and Applications\nThis application is an improvement over applicant's U.S. Pat. No. 3,479,838 (1969) entitled \"Cooling Material\". It is also a continuation-in-part of applicant's pending application Ser. No. 383,004 entitled \"Method and Apparatus for Monitoring and Diagnosing Peripheral Blood Flow\" and Ser. No. 673,509 entitled \"Blood Flow Monitoring Device\" now Pat. No. 4,569,335 (1986).\n2. Field of the Invention\nThis invention relates to fabric or garments for localized cooling of a person's skin or of another object.\n3. Need For the Invention\nLocalized cooling is extremely useful in many applications of which the following examples are but a few. Industrial workers in high temperature environments such as foundries and agricultural workers in the desert are exposed to high temperatures. At high temperatures, workers fatigue more easily and productivity decreases. It is impractical to air condition an entire foundry; farm laborers also cannot work in air conditioned environments. Therefore, it is desirable to cool these workers' skin. Peak efficiency and comfort may still be maintained if only the trunk or trunk and head are cooled. The limbs, which have a large surface area to volume ratio, naturally radiate enough heat. As a matter of fact, amputees often find it difficult to radiate sufficient heat from their bodies even at low levels of activity in environments that would be relatively comfortable for persons with all of their limbs. Thus, cooling of the trunk of amputees is another application of the present invention.\nSome persons are exposed to heat from one localized source. For example, a race car driver with a rear engine car is exposed to high heat quantities at his or her back. For racing efficiency, it is not practical to air condition the driver compartment, and added insulation adds undesirable weight to the vehicle. Moreover, race drivers usually wear protective clothing for fire and crash protection that tends to retain body heat. For them, localized cooling, which might concentrate on cooling the driver's back, would be most helpful.\nSurgical patients, especially those undergoing open heart surgery, often must be cooled to lower metabolism and blood flow. This is normally done with ice, but the ice needs replenishing, and it is more difficult to control the actual temperature.\nAs another example, one can cool a small area only and give the user the perception of comfort. For example, a cool head band or cap worn during strenuous exercise may allow the user to feel comfortable even though the rest of the body is above normal temperature.\nLocalized applications of cooling are not confined to humans. Electronics components, for example, heat up. Although they may be ventilated with outside air, as more circuitry is confined to a smaller area, outside air may be insufficient or at too high a temperature to cool the components. It would be useful, therefore, to provide a small device that could be inserted into an electronic device merely for the purpose of cooling localized areas. There may also be certain chemical processes which would benefit from localized cooling.\nLocalized cooling can be obtained by refrigerated air conditioning, but in many of the applications just discussed, that type of air conditioning is not practical because of weight, power consumption, restrictions on mobility and/or other reasons.\n4. Prior Art\nApplicant's Pat. No. 3,479,838 teaches a cooling fabric primarily for use in space. The cooling fabric had several layers, and water is injected between some of the layers. In the vacuum of outer space, water boils at skin temperature, and heat from the skin causes the water to boil. Heat exchange during a change of phase is great, and the resultant boiling caused the skin to cool. The device does not work at or near standard temperatures (body or environment) and pressure.\nThere are garments with channels for carrying cooling water. Applicant's earlier patents suggest the use of additional cooling liquid to augment the cooling from the boiling water. That type of system is prohibitive for everyday use on the ground because of weight and power requirements. Recirculating water must be cooled somehow unless one has access to cooling water that can be expelled after use or refrigerated at the expense of energy. Also, if the garment uses cold water or other liquid, the liquid flows through tubes, which must have thick enough walls to avoid kinking. These tubes have poor surface area contact and poor heat transfer characteristics.\nTo cool the skin while taking advantage of boiling at atmospheric pressure, one needs a liquid with a boiling point at atmospheric pressure below normal ambient temperatures. Heat from the skin boils these liquids. These refrigerant liquids are known and are used in air conditioning systems. Various fluorinated hydrocarbon fluids, such as those sold under the trademark Freon and mixtures of Freons, are known refrigerant liquids, and the present invention uses them. For cooling skin, the boiling point at standard pressure should be above 0.degree. C. to prevent freezing of the skin. Liquids boiling below 0.degree. C. could be used for cooling other objects.\nIn applicant's pending applications, localized cooling is used to cool a very small area of skin for a medical application. The cooling is accomplished by means of vaporization of a refrigerant liquid. The resultant gas is expelled to the atmosphere. The gas is costly and considered by some to pollute, and some have an offensive odor; it may be objectionable to vent the gas to the atmosphere, but these factors are minor in the medical device because it vents relatively small amounts of refrigerant gas. In a whole-body cooling apparatus, a much greater volume of gas would be lost, which raises the cost of each application. Therefore, it is an object of the present invention to disclose and to provide a localized cooling material in which the vaporized gas is not vented to the atmosphere.\nRecycling the gas creates a problem. One cannot allow the gas to stay in the cooling device or the device quickly becomes pressurized and boiling efficiency drops dramatically. It is impractical and inefficient to repressurize the relatively small amount of low pressure gas immediately and direct it back to the storage tank. It is an object of the present invention, therefore, to disclose and provide a system for holding interim storage of the spent gas and injecting it back into the storage tank as a liquid.\nFor maximum efficiency and control, the liquid should remain liquid until it is in a location where the skin temperature causes it to boil. Premature boiling in feed lines wastes the cooling potential of the liquid. Both of the pending applications attempt to cause the injection of the liquid to take place in the cooling chamber to conserve the cooling potential of the liquid so that it is not wasted in the lines. Also, by causing only liquid to be injected into the cooling chamber, control may be maintained. It is an object of the present invention, therefore, to disclose and provide valves or other controllers as close to the cooling location as possible. In the medical devices, very precise temperature controls to .+-.0.05.degree. F. (.+-.0.03.degree. C. [metric equivalents are approximate]), may be needed. A cooling fabric would not have to cool in such precise increments. A .+-.1.degree. or 2.degree. F. difference may be acceptable. Another object of this invention is to disclose and provide a relatively low-cost valve that can still be used to maintain the temperature control needed.\nIf the fabric covers wide areas of the body, it is desirable to be able to control those areas to effect different rates of cooling and different temperatures at localized sites. As a result, an additional object of the present invention is to design a valve that uses low power but is small and compact enough so that it can fit into the cooling fabric at one or several sites without being heavy, conspicuous and expensive.\nApplicant's application Ser. No. 673,509 discloses (in somewhat schematic form) a cooling fabric which uses injectors disclosed in that application. It is an object of the present invention to improve on that system and utilize the fabric as a cooling chamber and to improve on the injectors themselves. These and other objects of the present invention will become evident from the following description of it."}
{"text": "The existing display screen usually includes a display screen frame and a display screen installed in the display screen frame, wherein a signal connector of the display screen is connected to a signal interface on the display screen frame. Since the connection between the signal connector of the display screen and the signal interface on the display screen frame has certain tightness, and both the frame and the display screen are provided with a magnet, and the two are adsorbed by the magnets, so that the display screen and the frame are tightly installed. Therefore, when a display module needs to be unloaded from a display screen box, the display screen needs to be adsorbed by a magnetic auxiliary screen disassembly tool, so as to disassemble the display screen by offsetting the magnet adsorption force between the display screen and the frame. However, this situation is prone to have the problem that the lamp beads of the display screen are damaged by the magnetic auxiliary screen disassembling tool; moreover, on one hand, the magnetic auxiliary screen disassembling tool has a heavier weight, after the screen is disassembled, the magnetic auxiliary screen disassembling tool is possible to fall off, and the magnetic auxiliary screen disassembly tool is much heavier and is not convenient to operate."}
{"text": "The present invention relates to liquid crystal display devices, and, more particularly, to liquid crystal display devices of the lateral electric field type.\nA liquid crystal display device of the lateral electric field-driven type comprises a pair of transparent substrates disposed so as to be opposed to each other and to be spaced by a liquid crystal (LC) layer sealed therebetween, and a pixel electrode and a counter electrode are spaced apart from the pixel electrode formed in each of the pixel regions on a surface at the liquid crystal layer side of one of the pair of transparent substrates. A pixel electrode and a counter electrode disposed at each of the pixel regions generates an electric field therebetween so as to control the optical transmissivity of the liquid crystals in the vicinity of the pixel region. Since the liquid crystal display device of the lateral electric field type generates an electric field having a major component extending substantially parallel to a main surface of the transparent substrate between the pixel electrode and the counter electrode for switching the orientations of the liquid crystal molecules, the lateral electric field scheme is also referred to as the xe2x80x9cIn-Plane-Switching (IPS) schemexe2x80x9d also.\nAn active-matrix liquid crystal display device employing the lateral electric field scheme has a plurality of scan signal lines and a plurality of counter voltage signal lines both of which extend in an xe2x80x9cxxe2x80x9d direction (row direction), while being juxtaposed in parallel in a xe2x80x9cyxe2x80x9d (column) direction (transverse to the x direction), and a plurality of video signal lines which extend in the y direction, while being juxtaposed in the x direction on the surface at the liquid crystal layer side of the one of the pair of transparent substrates. When the scan signal lines, the counter voltage signal lines and the video signal lines are disposed on a surface of one of the pair of transparent substrates, each of the pixel regions is defined as an area surrounded by one of the scan signal lines, one of the counter voltage signal lines disposed adjacent to the scan signal line, and a pair of the image signal lines disposed adjacent to one another. Each of the pixel regions has a switching element, a pixel electrode, and a counter electrode. The switching element operates so as to be turned on by a scanning signal supplied by one of the scanning signal lines. The switching element, when turned on, supplies a video signal being transmitted by one of the video signal lines thereto to the pixel electrode in a pixel region corresponding thereto. A counter voltage for generating a potential difference between the counter electrode and the pixel electrode in the pixel region is transmitted through one of the counter voltage signal lines and is applied to the counter electrode.\nA liquid crystal display device having the structure described above is characterized by the fact that it produces an image having a higher contrast ratio and exhibits a remarkably wider viewing angle.\nThe lateral electric field-type liquid crystal display device referred to above is described in detail in such documents as Japanese Patent Application Laid Open Hei 05-505247 (No. 505247/1993), Japanese Patent Publication Sho 63-021907 (No. 021907/1988), and Japanese Patent Application Laid Open Hei 06-160878 (No. 160878/1994).\nThe lateral electric field-type liquid crystal display device having the structural features described above utilizes electrodes (e.g. an pixel electrode and a counter electrode) formed of thin conductive films of several thousands Angstrom (xc3x85) (thin film electrodes, hereinafter) and generates an electric field substantially in parallel with a surface of one of the pair of transparent substrates. However, it is difficult to reproduce the aforementioned electric field sufficiently and accurately in the liquid crystal layer lying between the thin film electrodes.\nIn addition to the aforementioned difficulty, the aforementioned electric field being generated between the pixel electrode and the counter electrode in a pixel corresponding thereto is affected by another electric field appearing from the video signal line adjacent to the pixel.\nOn the other hand, a black matrix layer (a light shielding film) is formed on a surface of another of the pair of transparent substrates confronting the one of the pair of transparent substrates on which the pixel electrode and the counter electrode are formed. For instance, the black matrix layer is formed on the surface of the substrate which confronts the liquid crystal layer. The black matrix layer may be formed of an insulating material or the like having such a high resistivity that the electric field between the pixel electrode and the counter electrode is generated to be as strong as possible.\nThe black matrix, being formed of an insulating material or the like, tends to have a higher optical transmissivity than that of metals so that a problem of deterioration of the contrast ratio occurs in an image displaying operation of the liquid crystal display device.\nThe present invention has been made in view of the technical background mentioned above, and an object of the present invention is to provide an improved liquid crystal display device which is suitable for generating an electric field sufficiently between the pixel electrode and the counter electrode as to control the orientations of the liquid crystal molecules lying therebetween in spite of utilizing a simple structure.\nAnother object of the present invention is to improve the contrast ratio of an image displayed by the liquid crystal display device.\nSome representative aspects and features of the present invention as disclosed herein will be briefly summarized as follows.\nThe present invention provides a liquid crystal display device comprising: a pair of substrates being superposed on one another with a predetermined space therebetween in which a liquid crystal layer is sealed; at least one video signal line being formed on a surface of one of the pair of substrates confronting the liquid crystal layer; and at least one pixel being formed on the surface of the one of the pair of substrates and having at least one pixel electrode to which a signal is supplied through the at least one video signal line and at least two counter electrodes being spaced from the at least one pixel electrode, one of the at least two counter electrodes being disposed adjacent to the at least one video signal line. The display device is characterized in that: (1) the at least one video signal line and the one of the at least two counter electrodes are spaced from one another by at least two insulating regions disposed along the surface, and (2) one of the at least two insulating regions comprises a material having a higher permittivity than that of the other of the at least two insulating regions.\nThe present invention also provides another liquid crystal display device comprising: a pair of substrates being disposed opposite to one another a liquid crystal layer being sealed therebetween; at least one video signal line formed on a surface of one of the pair of substrates confronting the liquid crystal layer; and at least one pixel being formed on the surface of the one of the pair of substrates, being arranged adjacent to the at least one video signal line, and having at least one pixel electrode and at least one counter electrode, so as to apply an electric field to the liquid crystal layer therebetween. The display device is characterized in that: (1) at least a surface of the at least one video signal line is covered with an insulating material having a permittivity (a dielectric constant) higher than 6.7.\nThese and other objects, features and advantages of the present invention will become more apparent from the following description when taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device including a circuit comprising a thin film transistor (hereinafter referred to as a TFT) over a substrate having an insulating surface and a method of fabricating the same. Particularly, the invention relates to a structure of an electro-optical device typified by a liquid crystal display device, and an electronic equipment incorporating the electro-optical device, or relates to a structure of an electro-optical device typified by an EL (electro luminescence) display device using an EL material capable of obtaining electro Luminescence, and an electronic equipment incorporating the electro-optical device.\nIncidentally, in the present specification, the term “semiconductor device” indicates any devices functioning by using semiconductor characteristics, and includes the foregoing electro-optical device and the electronic equipment incorporating the electro-optical device in its category.\n2. Description of the Related Art\nA semiconductor device including a large area integrated circuit made of TFTs is under development. An active matrix type liquid crystal display device, an EL display device, and a contact type image sensor are its typical examples.\nThe TFT can be classified according to its structure and fabricating method. Particularly, since a TFT (crystalline TFT) including a semiconductor film having crystal structure as an active layer has a high field effect mobility, it has been possible to form various functional circuits.\nIn the present specification, the semiconductor film having the crystal structure includes a single crystal semiconductor, a polycrystal semiconductor, and a microcrystal semiconductor, and further, includes a semiconductor disclosed in Japanese Patent Application Laid-open No. Hei. 7-130652, No. Hei. 8-78329, No. Hei. 10-135468, No. Hei. 10-247735, or No. Hei. 10-135469. The disclosure of the above Japanese Patent Applications is incorporated with herein by reference.\nIn an active matrix type liquid crystal display device, for every functional block, a pixel region (also called a pixel matrix circuit) comprising n-channel TFTs and a driver circuit including a CMOS circuit as a basic unit, such as a shift register circuit, a level shifter circuit, a buffer circuit, and a sampling circuit, are formed over one substrate.\nIn the contact type image sensor, an integrated circuit such as a sample-and-hold circuit, a shift register circuit, and a multiplexer circuit is formed by using TFTs.\nSince these circuits do not necessarily have the same operation condition, characteristics required for TFTs have been naturally different not a little from one another.\nThe characteristics of a field effect transistor such as a TFT can be considered by dividing them into a linear region where a drain current increases in proportion to a drain voltage, a saturation region where even if a drain voltage is increased, a drain current is saturated, and a cut-off region where even if a drain voltage is applied, a current does not flow ideally. In the present specification, the linear region and the saturation region are referred to as an ON region of a TFT, and the cut-off region is referred to as an OFF region. For convenience, a drain current in the ON region is referred to as an ON current, and a current in the OFF region is referred to as an OFF current.\nA pixel portion comprises a switching element made of an n-channel TFT (hereinafter referred to as a pixel TFT) and an auxiliary holding capacitance, and applies a voltage to a liquid crystal to drive it. Here, the liquid crystal is required to be driven by an alternating current, and a system called frame inversion driving has been adopted. Thus, as characteristics of TFTs to be required, it has been necessary that the OFF current is sufficiently reduced.\nSince a buffer circuit of a driver circuit is applied with a high driving voltage, it has been necessary to increase withstand voltage. Besides, in order to increase current driving performance, it has been necessary to sufficiently secure the ON current.\nHowever, there has been a problem that the OFF current of a crystalline TFT is apt to become high. The crystalline TFT has been regarded as being inferior to a MOS transistor (transistor fabricated on a single crystal semiconductor substrate) used for an LSI or the like in reliability. For example, a deterioration phenomenon such as a lowering of ON current has been sometimes observed in the crystalline TFT. It has been considered that this cause is a hot carrier effect, and hot carriers generated by a high electric field in the vicinity of a drain cause the deterioration phenomenon.\nAs a structure of a TFT, a low concentration drain (LDD: Lightly Doped Drain) structure has been known. In this structure, an impurity region having a low concentration is provided between a channel forming region and a source region or drain region added with an impurity at a high concentration, and this low concentration impurity region is called an LDD region.\nAccording to positional relation to a gate electrode, the LDD structure includes a GOLD (Gate-drain Overlapped LDD) structure where it overlaps with the gate electrode, an LDD structure where it does not overlap with the gate electrode, and the like. The GOLD structure has been able to relieve a high electric field in the vicinity of a drain, to prevent the hot carrier effect, and to improve the reliability. For example, in “Mutsuko Hatano, Hajime Akimoto and Takeshi Sakai, IDEM 97 TECHNICAL DIGEST, p 523-526, 1997”, it is ascertained that extremely excellent reliability can be obtained in the GOLD structure of a side wall comprising silicon as compared with TFTs of other structures.\nOn the other hand, as another problem in relation to a large area integrated circuit, there has been a problem of wiring. An integrated circuit comprising TFTs is provided with a gate wiring line connected to a gate electrode and a data wiring line connected to a source electrode or drain electrode. Particularly, the gate wiring line has a problem of a wiring delay due to influence of parasitic capacitance and wiring resistance. Although a material such as molybdenum (Mo), tantalum (Ta), or tungsten (W) has been used for the gate electrode and the gate wiring line in view of heat resistance, these have a sheet resistivity of about 10 Ω, and have not been suitable for a large area integrated circuit. It has been originally preferable to use a low resistance material such as aluminum (Al) or copper (Cu).\nHowever, the GOLD structure has a problem that the OFF current becomes high as compared with a normal LDD structure, and it has not been necessarily preferable to form all TFTs with the GOLD structure in a large area integrated circuit. For example, in a pixel TFT, if the OFF current is increased, a power consumption is increased and an abnormality appears on image display. Thus, it has not been preferable to apply a crystalline TFT of the GOLD structure as it is.\nMoreover, the LDD structure has a problem that the ON current is decreased by an increase of series resistance. Although the ON current can be freely designed through a channel width of a TFT and the like, for example, it has not been always necessary to provide the LDD structure in the TFT constituting a buffer circuit."}
{"text": "Workability of concrete refers to the effort (or energy) required to manipulate a freshly mixed quantity of concrete with minimum loss of homogeneity. The manipulation may refer to pumping, placing, consolidation and/or finishing of the concrete. Workability of concrete is quantified in terms of “slump” or “slump flow”. The slump flow test has been standardized as ASTM C 1611, “Slump Flow of Self-Consolidating Concrete.” Slump and slump flow are measurements of concrete rheology and is determined using a slump cone. A slump cone is standardized measurement cone having a predefined volume and angle. FIG. 1A illustrates a standard slump cone 100 which includes a top opening 102 and a bottom opening 104.\nThe slump cone 100 is used by placing the slump cone 100 on a flat surface and filling the slump cone 100 with fresh concrete through the top opening 102 as shown in FIG. 1B. The slump cone 100 is completely filled and any excess concrete at the top of the slump cone 100 is scraped off. The slump cone 100 is lifted directly upwards. The concrete 110 slowly spreads out over the flat surface without the support of the slump cone 100 to hold the concrete 110 in place. The spreading action causes the height of the concrete 110 to decrease to a height 112 from an initial height 116 corresponding to the height of the slump cone 100. The distance 114 represented by the change in height of the concrete 110 is referred to as the “slump”. The “slump” is indicative of the yield value of the concrete. The concrete 110 ultimately spreads out over the flat surface to form a crude disc or pancake shape as shown in FIG. 1C. The diameter 118 of the disc formed by the concrete 110 is referred to as the “slump flow”.\nAnother measurement of the workability of the concrete is the T50 test. The T50 test quantifies the viscosity of the concrete 110 and measures the amount of time in seconds for the concrete in the slump flow test to spread to a diameter of 50 cm or 500 mm). The T50 test has been standardized as ASTM C 1611. The T50 time is indicative of the viscosity of the concrete. The slump, slump flow and T50 are used to assess the rheological properties of the concrete and predict how it will flow or move under the force of gravity or positive force into a desired shape or position.\nConventional concrete has a slump flow of 80 mm to 120 mm. The placement and finishing of concrete floors in residential and commercial applications using a conventional concrete is difficult, labour intensive and time consuming due to this low slump. Conventional concrete having a minimum compressive strength of 25 MPa is typically designed to have a maximum slump flow of 100 mm; however, water is added in the field to raise the slump flow to 200 mm to improve the workability of the concrete, and in particular the placeability of the concrete. The consequences of adding extra water to increase the water/cement ratio, is to decrease ultimate strength and increase shrinkage.\nSCC is a type of concrete characterized by a low yield, moderate viscosity which can be used to ensure a uniform suspension of aggregate during transportation, placement and finishing until the concrete sets. SCC has properties which are desirable in the construction of the concrete floors in residential and commercial applications; however, known SCC compositions are more costly to produce than conventional concrete and provide a surface which is prone to tearing and peeling in response to mechanical working/finishing (smoothing and/or hardening) such as when power troweled. This results in surface defects such as blisters and delaminations which are unacceptable for most applications.\nAccordingly, a SCC mixture is desired which is comparable in cost to conventional concrete used in residential and industrial floor construction but which requires less time to place and finish, thus increasing the productivity of the concrete supplier and user."}
{"text": "Of all of the applications where computer vision is used, face detection presents an extremely difficult challenge. For example, in images acquired by surveillance cameras, the lighting of a scene is usually poor and uncontrollable, and the cameras are of low quality and usually distant from potentially important parts of the scene. Significant events are unpredictable. Often, a significant event is people entering a scene. People are typically identified by their faces. The location and orientation of the faces in the scene are usually not controlled. In other words, the images to be analyzed are substantially unconstrained.\nIn applications where a large number of images are to be analyzed, or the analysis needs to be performed in real-time, the number of images that actually include faces is usually very small. Moreover, in most images that do include a face, the amount of the image that is a face is also very small. Therefore, it is desirable to rapidly eliminate images or portions of images that do not include a face, thus more processing can be devoted to portions of images that potentially contain a face.\nFace detection has a long and rich history. Some techniques use neural network systems, see Rowley et al., “Neural network-based face detection,” IEEE Patt. Anal. Mach. Intell., Vol. 20, pages 22–38, 1998. Others use Bayesian statistical models, see Schneiderman et al., “A statistical method for 3D object detection applied to faces and cars,” Computer Vision and Pattern Recognition, 2000. While neural network systems are fast and work well, Bayesian systems have better detection rates at the expense of longer processing time.\nU.S. Pat. No. 5,710,833, “Detection, recognition and coding of complex objects using probabilistic eigenspace analysis” issued to Moghaddam, et al. on Jan. 20, 1998, describes a system for detecting instances of a selected object or object feature, e.g., faces, in a digitally represented scene. Their method utilizes analysis of probability densities to determine whether an input image, or portion thereof, represents such an instance.\nU.S. Pat. No. 6,337,927, “Approximated invariant method for pattern detection” issued to Elad et al. on Jan. 8, 2002, describes a system and method for classifying input feature vectors into one of two classes, a target class (faces) and a non-target class (non-faces). The system utilizes iterative rejection stages to first label the input vectors that belong in the non-target class in order to identify the remaining non-labeled input vectors as belonging in the target class. The operation of their system is divided into an off-line (training) procedure and an online (actual classification) procedure.\nDuring the off-line procedure, projection vectors and their corresponding threshold values that will be used during the on-line procedure are computed using a training set of sample non-targets and sample targets. Each projection vector facilitates identification of a large portion of the sample non-targets as belonging in the non-target class for a given set of sample targets and sample non-targets.\nDuring the on-line procedure, an input vector is successively projected onto each computed projection vector and compared with a pair of corresponding threshold values to determine whether the input vector is a non-target. If the input vector is not determined to be a non-target during the successive projection and thresholding, the input vector is classified as a target.\nAs a disadvantage, the features used by their method are limited to linear projections and upper and lower thresholds. Thus, outcomes of early features do not effect the outcomes of later tests. In other words, their tests can be reordered with no effect on the final result. In addition, because the classifier is a conjunction of linear projections followed by a threshold, the resulting classification function is limited to a convex region in the input space, which is a strong limitation that makes their system a weak classifier.\nA new framework for feature detection is described by Viola et al., in “Rapid Object Detection using a Boosted Cascade of Simple Features,” Proceedings IEEE Conf. on Computer Vision and Pattern Recognition, 2001. They present three new insights: a set of image features which are both extremely efficient and effective for face analysis, a feature selection process based on Adaboost, and a cascaded architecture for learning and detection. Adaboost provides an effective learning algorithm and strong bounds on generalized performance, see Freund et al., “A decision-theoretic generalization of on-line learning and an application to boosting,” Computational Learning Theory, Eurocolt '95, pages 23–37. Springer-Verlag, 1995, Schapire et al., “Boosting the margin: A new explanation for the effectiveness of voting methods,” Proceedings of the Fourteenth International Conference on Machine Learning, 1997, Tieu et al., “Boosting image retrieval,” International Conference on Computer Vision, 2000.\nHowever, there are several problems with their cascaded approach. The simplest is that the cascade allows only for rejection thresholds, and does not have acceptance thresholds. This degrades performance when the distribution of acceptable and rejected examples are substantially equal.\nAnother disadvantage of their method is that later stages of the cascade ignore detailed information that was available during earlier stages of the cascade. If the early cascade stages accept features only marginally close to the threshold, then this useful information plays no role in later decisions to reject or accept the feature. For example, it is possible for multiple stages to accept a given patch because it is significantly above the rejection threshold, and then subsequently reject that patch as soon as it falls below the rejection threshold in later stages. In this case it is often better to accept the given patch.\nAnother disadvantage is their mechanism for setting the rejection thresholds. Their rejection thresholds are set for each cascade stage starting from the first stage. For any given task, their goal is to detect most of the objects, though it is impossible to detect all objects of a particular class. In the case of face detection, perhaps 90% to 95% of the faces are detected. The rejection threshold is set so that a very high percentage of the examples are detected. Therefore, in the early stages, one must be very conservative and not discard a potentially correct patch. For early stages in the cascade, 99% or more of the objects must be detected in order to ensure that valuable patches are not discarded.\nTheir process does not distinguish example patches that are easy to classify, and are eventually classified correctly, and those example patches, which are very difficult to classify and are eventually misclassified. If one could determine which example patches would eventually be misclassified, after the entire cascade, then those example patches could be eliminated early with no loss to the accuracy of the overall process. Because the rejection threshold setting process starts with the early stages of the cascade, these “difficult” examples cannot be identified so the rejection threshold is set for all positive examples. This results in rejection thresholds which are lower and less effective in reducing processing time.\nTherefore, there is a need for an object detection system that improves upon the prior art."}
{"text": "This invention relates to electric switches of the type which have an opposed pair of stationary contacts and a movable contact and are fabricated by pressing, to an apparatus for the manufacture thereof and to a process for said manufacture.\nThe quality and performance of electric switches have no little effect upon the quality of various electric appliances which these switches serve. In an extreme case, failure in an electric switch can even endanger the safety of a person handling the electric appliance containing that switch. All the electric switches are, therefore, required to possess stable quality, provide quick operation and preclude troubles. For this reason, ordinary electric switches currently manufactured require much time and labor in jobs of inspection and possible readjustment which are performed subsequent to conclusion of the assembly of parts. U.S. Pat. No. 3,189,703 has issued to an invention which relates to a method for providing epochal improvements in the manufacture of electric switches. It provides a method which comprises the steps of punching sets each of three terminals in an interconnected form out of a strip of electroconductive material, mounting on each set an opposed pair of stationary contacts and a leaf spring fitted with a movable contact, embedding them with a molded plastic webbing and thereafter removing the webbing temporarily interconnecting the three terminals. This invention is excellent in that there is no possibility of dimensional error because the three terminals and the conductor are always retained in their relative position acquired at the time of punching and that the switches are suitable for mass-production. Nevertheless, the work of mounting a leaf spring and an opposed pair of contacts on such terminals which are manufactured efficiently by the method of said U.S. patent is not efficient. Particularly, the stationary contacts are formed by a troublesome step of mounting two prefabricated mushroom-shaped silver contacts in the slots and peening the divergent portion of the contact at predetermined positions into the slot. Contacts of the type formed by pressing silver rods into hemicircular slots are not totally reliable in terms of joint strength and electroconductivity. Worse yet, the very low dimensional accuracy of contact gap poses itself as a serious problem.\nAccording to said U.S. patent, each unit of interconnected terminals is embedded by molded plastic webbing. Inclusion of such troublesome and inefficient molding work in the process of switch production can hardly be regarded as an improvement.\nThe main object of the present invention is to provide electric switches which are suitable for mass-production and which enjoy so precise fabrication as to render the jobs of inspection and readjustment practically unnecessary.\nAnother object of this invention is to provide electric switches each possessed of an opposed pair of contacts, completed by a sole operation of pressing using a strip of contact material such as of silver which is fed to the electric switch production line.\nStill another object of this invention is to provide a process for completing high-quality stationary contacts in electric switches by a sole operation of pressing using a strip of contact material such as of silver.\nYet another object of this invention is to provide electric switches which are completed by a sole operation of pressing using a strip of contact material such as of silver and subsequently enclosed in prefabricated cases with high efficiency.\nAnother object of this invention is to provide electric switches which, on account of a leaf spring fitted at one end thereof with a movable contact and mounted on the other end thereof around the neck of a common inner end of terminals, enjoy sharp snap action and durability and which, at the same time, enjoy freedom from a conventional disadvantage that the snap action is adversely affected by the degree with which the spring is fastened."}
{"text": "The subject matter hereof is related to the subject matter of U.S. Pat. No. 4,061,389 and is directed specifically to the improvement that promotes effective connection of a pull tool to a down hole object or fish and simple and effective disconnection of the pull tool from the article if such is desired.\nDuring drilling operations, well servicing operations and the like, objects may become lost or stuck within a well bore and these objects are typically referred to in the industry as \"fish.\" For example, a section of drill pipe or production tubing might become disconnected from a pipe string and it is then necessary to introduce a pull tool into the well bore, establish connection with the fish and then remove the fish. Many different types of pull tools, grap pling devices, spears, etc. have been developed for the purpose of conducting fishing operations in well bores. It is typical for wire line operations to be utilized for the purpose of introducing a pull tool into the well bore for fishing operations.\nIn many cases, down hole devices are positioned within casing and tubing strings for the purpose of controlling well operations. In most cases, such down hole devices are equipped with API Standard fishing necks which normally include an undercut shoulder to enable grappling by standard installation and retrieval tools. Quite often the API Standard fishing neck will become corroded or eroded by well conditions to the point that only a stub pipe is exposed without the usual undercut shoulder that is provided on the fishing neck. Connection between a wire line controlled fishing tool and a worn fishing neck may be accomplished by means of an overshot type grappling device generally defined by a collet structure having internal teeth that establish a gripping relation with the worn fishing neck. Although overshot type retrieving tools are successfully utilized in many cases, a common problem with such tools is the inability of the operator to achieve disconnection from the worn fishing neck in the event the tool is unable to accomplish an effective pulling operation. Occasionally, disconnection can be achieved by a substantial jarring or other violent mechanical movement, but, in some cases, disconnection of the pulling tool from the fish is extremely difficult if not impossible. In cases where the pulling tool or wire line is inadequate for the pulling operations that are required, it is necessary that the pulling tool be disconnected from the fish and replaced with a wire line controlled pulling tool of substantially greater pulling capacity. If a light-weight wire line pulling tool is unable to accomplish the pulling operation and becomes firmly fixed to the fish or other object to be pulled, retrieval of the fishing equipment itself obviously compounds the problem and adds materially to the expense of the service operation.\nAs mentioned above, application of shocks to the pulling tool will sometimes result in disconnection, thereby enabling service personnel to remove the wire line pulling tool and substitute a tool and wire line of substantially heavier gauge. Application of mechanical shocks to the down hole equipment of the well and the wire line tool itself can easily cause failure or excessive wear of one or more of these mechanical components. It is desirable, therefore, to provide down hole service equipment, such as wire line controlled pulling tools, that will effectively become interconnected with down hole objects such as wire line tools, fish, etc., and, in the event disconnection is necessary, will readily become disconnected upon simple mechanical movement of a wire line controlled pulling tool."}
{"text": "The present invention generally relates to rivets and more particularly to a blind rivet assembly.\nIt is known to set blind rivets with a mandrel. For example, reference should be made to U.S. Pat. No. 5,689,873 entitled “Tacking Fastener” which issued to Luhm on Nov. 25, 1997, Great Britain Patent No. 286,471 entitled “Improvements in and Relating to Rivets and Riveting,” and Great Britain Patent No. 2 150 661 entitled “Pull-Type Blind-Riveting Assemblies.” These rivets, however, outwardly extend past outside surfaces of workpieces after rivet setting. This disadvantageously takes up too much space, and is prone to interference or damage during assembly. More recently, flush blind rivets, using breakable mandrels, have been employed. While these constructions have significantly improved the industry, further improvements are desirable."}
{"text": "Current sports shoes for athletes have a plastic sole and textile upper. Such shoes are manufactured to be as light as possible in order to improve the performance of the athletes, particularly over short distances, e.g. 100 m or 200 m.\nOne example of a sports shoe is known from European Patent Application No. EP-A-0 799 583 (Technica S.p.A). The shoes described in this patent have a sole made from a resilient material such as rubber with a strengthening insert. This shoe is not suitable for short distance track events because of its weight.\nA variety of athletics shoes are produced by several manufacturers. The Saucony middle distance race shoe has a very light-weight fabric to keep the total weight of the shoe as low as possible. An Asics sprinting spike shoe uses a rubber grip on the sole of the shoe from the footplate backwards.\nSandals are also known in the art which can be used for trail running and hiking. These sandals, such as those manufactured by the Teva Company, Flagstaff, Ariz., United States, have an open structure with straps to hold the sandals on the foot. The sandals have a substantially flat sole with an open toe, arch and heel areas. Whilst such sandals, such as the Wraptor 2 model, may be suitable for long-distance trail running their heavy weight makes them unsuitable for use by track athletes.\nCycling shoes are known in the art which are designed to be as light as possible and also to allow as much transmission of the athlete's muscular power to the cycle. For example, U.S. Pat. No. 5,836,094 (Figel) teaches a cycling shoe in which a composite made of fibre and resin replaces the sole and part of the upper of a traditional shoe. In this patent, the complete sole of the shoe is made from the composite. Manufacture of the composite is from carbon fibre, Kevlar fibres and epoxy resin which are built up using layers of carbon fibre sheets and Kevlar layers.\nAnother cycling shoe is described in U.S. Pat. No. 5,406,723 (Okajima, assigned to Shimano, Inc., Osaka, Japan). This patent describes a sole of a shoe with a fibrous layer surrounding a core made of polyurethane foam. A glass cloth is applied to the upper and lower surfaces of the core to form a layer. A carbon fibre sheet impregnated with phenol resin is then placed over the glass cloth followed by an impregnated glass cloth. This process is completed until a stiff sole is manufactured.\nCycling shoes are, however, not suitable for use by track athletes because they have a stiff sole which is substantially inflexible. Track athletes require a shoe which is flexible in order to maximise their performance, particularly at the starting blocks."}
{"text": "1. Field of the Invention\nThe present invention relates to nonvolatile memory. More particularly, the present invention relates to methods and apparatus for injecting electrons and holes onto the floating gate of a floating gate MOS transistor nonvolatile memory cell.\n2. The Background\nNonvolatile memories employing floating gate technology have found wide use in a variety of applications, and have become an increasingly important type of semiconductor nonvolatile memory. In a floating gate nonvolatile memory, the floating gate electrode of a MOS transistor is electrically isolated from neighboring electrodes by a high quality dielectric that surrounds the floating gate. The contents of the memory, e.g. whether a xe2x80x980xe2x80x99 or a xe2x80x981xe2x80x99 in a digital memory, is determined by the amount of charge on the floating gate.\nIn a memory known to those of ordinary skill in the art as flash memory, electrons are commonly added to the floating gate by a process of hot electron injection that is controlled by the bias applied to various elements of the floating gate transistor. To remove electrons from the floating gate, electrons tunnel from the floating gate to surrounding electrodes under the influence of a high electric field. In the most common form of this technology, the electrons are caused to tunnel from the floating gate to an underlying silicon region by applying a bias across a relatively thin layer of silicon dioxide. In order that the charge be retained on the floating gate for an extended period of time, it is important that the tunneling operate in what is referred to by those of ordinary skill in the art as the Fowler-Nordheim mode so that the current under cell storage and read conditions is very low, typically less than 10xe2x88x9215 A/cm2.\nIt has been demonstrated that to restrict tunneling to the Fowler-Nordheim mode, the oxide through which the electrons tunnel should be thicker than about 5 nanometers. Further, the 5 nm lower bound is not a practical limit because the act of applying the voltage across the oxide to cause tunneling damages the oxide. To reduce the stress induced leakage current that occurs due to the damage to the oxide, the minimum oxide thickness should be increased from the lower bound of about 5 nm to about 8 nm. Naruke, K., et al, xe2x80x9cStress Induced Leakage Current Limiting to Scale Down EEPROM Tunnel Oxide Thicknessxe2x80x9d, IEDM Tech. Digest pp. 424-7 (1988).\nAs integrated circuit (IC) dimensions have been scaled below 0.25 xcexcm, the operating voltages of CMOS circuits have also been scaled down. Unfortunately, it has been demonstrated that the applied voltages typically employed to inject electrons onto and tunnel off of the floating gate of an electrically erasable programmable read only memory (EEPROM) are not capable of being scaled down to the same degree as the operating voltages of below 0.25 xcexcm processes. Yoshikawa, K. et al, xe2x80x9cFlash EEPROM Cell Scaling Based on Tunnel Oxide Thinning. Limitationsxe2x80x9d, Symp. on VLSI Technol. Dig. Tech. Papers, pp. 79-80 (1991); Yamaguchi, Yoshiko, et al, xe2x80x9cONO Interpoly Dielectric Scaling Limit for Non-volatile Memory Devicesxe2x80x9d, Symp. on VLSI Technol. Dig. Tech. Papers, pp. 85-86 (1993); Caywood, J. M. and Gary Derbenwich xe2x80x9cNonvolatile Memoriesxe2x80x9d, in ULSI Device Technology, ed. S. M. Sze and C. Y. Chang, in press.\nBecause charge leakage through the tunnel oxide limits the scaling of the tunnel oxide thickness and the physics of the phenomena used for injecting charge onto the floating gate limit the scaling of the voltages, it is well understood in the art that the mismatch between the minimum physical dimensions and minimum operating voltages of the floating gate nonvolatile memories and the surrounding CMOS logic technology is becoming increasingly acute. Accordingly, what is needed are new apparatus and methods for charging and discharging the floating gate so that the physical dimensions and required voltages can be reduced to values more compatible with CMOS logic technology.\nA tunneling charge injector that includes a conducting injector electrode, a grid insulator disposed adjacent to the conducting injector electrode, a grid electrode disposed adjacent to the grid insulator, a retention insulator disposed adjacent to the grid electrode, and a floating gate disposed adjacent to the retention insulator, may be employed to inject charge from the conducting injector electrode onto a floating gate.\nIn one aspect of the present invention, the tunneling charge injector may be employed to inject both electrons and holes onto the floating gate electrode. Electrons are injected onto the floating gate when the conducting injector electrode is sufficiently negatively biased with respect to the grid electrode and the floating gate is biased to collect electrons from the grid. Holes are injected onto the floating gate when the conducting injector electrode is sufficiently positively biased with respect to the grid electrode and the floating gate is biased to collect, holes from the grid.\nIn another aspect of the present invention, separate tunneling electron and hole injectors may be employed to inject electrons or holes onto the floating gate. Electrons are injected onto the floating gate when the conducting injector electrode is sufficiently negatively biased with respect to the grid electrode and the floating gate is biased to collect electrons from the grid. Holes are injected onto the floating gate when the conducting injector electrode is sufficiently positively biased with respect to the grid electrode and the floating gate is biased to collect holes from grid.\nIn another aspect of the present invention, the tunneling charge injector may be employed in a nonvolatile memory cell having a nonvolatile memory element with a floating gate such as a floating gate MOS transistor. In the nonvolatile memory cell, the floating gate of the tunneling charge injector is coupled to or forms a part of the floating gate of the nonvolatile memory element. The tunneling charge injector is employed to inject charge onto the floating gate of the nonvolatile memory element.\nIn another aspect of the present invention, separate tunneling electron and hole injectors may be employed in a nonvolatile memory cell having a nonvolatile memory element with a floating gate such as a floating gate MOS transistor. In the nonvolatile memory cell, the floating gates of the separate tunneling electron and hole injectors are coupled to or form a part of the floating gate of the nonvolatile memory element. The tunneling electron injector is employed to inject electrons onto the floating gate of the nonvolatile memory element, and the tunneling hole injector is employed to inject holes onto the floating gate of the nonvolatile memory element.\nIn another aspect of the present invention, a memory device includes an array of nonvolatile memory cells wherein each of the memory cells comprises a nonvolatile memory element with a floating gate such as a floating gate MOS transistor and a tunneling charge injector having a floating gate that is either coupled to the floating gate of the nonvolatile memory element or forms a portion of the floating gate of the nonvolatile memory element.\nIn another aspect of the present invention, a memory device includes an array of nonvolatile memory cells wherein each of the memory cells comprises a nonvolatile memory element with a floating gate such as a floating gate MOS transistor, a tunneling electron injector having a floating gate that is either coupled to the floating gate of the nonvolatile memory element or forms a portion of the floating gate of the nonvolatile memory element, and a tunneling hole injector having a floating gate that is either coupled to the floating gate of the nonvolatile memory element or forms a portion of the floating gate of the nonvolatile memory element.\nIn another aspect of the present invention the tunneling charge injectors and memory devices described herein are equipped with a retention insulator featuring a first band gap in a first portion of said retention insulator disposed adjacent to said grid electrode and a second, greater band gap in a second portion of said retention insulator disposed between said grid electrode and said floating gate to facilitate charge injection."}
{"text": "This invention relates to static random access memory (xe2x80x9cSRAMxe2x80x9d) devices, and, more particularly, to a system and method for powering-up SRAM devices having PMOS access transistors to limit the initial current draw of such SRAM devices.\nMany integrated circuit devices, such as microprocessors, include on-board memory devices, such as SRAM devices. For example, SRAM devices are commonly used as cache memory because of their relatively fast speed. SRAM devices are also sold as stand-alone integrated circuits for use as cache memory and for other uses. SRAM devices are also more suitable for use as cache memory than dynamic random access memory (xe2x80x9cDRAMxe2x80x9d) devices because they need not be refreshed, thus making all SRAM memory cells continuously available for a memory access.\nFIG. 1 is a block diagram of a portion of a typical array 10 of SRAM cells 12 arranged in rows and columns. A plurality of complementary digit line pairs D, D* are used to couple complementary data to and from the memory cells 12 in a respective column. Several digit line pairs, typically 16 or 32 digit line pairs, are coupled to respective inputs of a column multiplexer 13. The column multiplexer 13 couples one pair of digit lines corresponding to a column address to a sense amplifier 14 and a write driver 16. The sense amplifer 14 provides a data output (not shown) indicative of the polarity of one digit line D relative to the other D* responsive to data being read from a memory cell 12 coupled to the selected digit line pair D, D*. The write driver 16 drives a differential voltage onto the digit lines D, D* to which the write driver 16 is coupled by the column multiplexer 13. The differential voltage applied between the digit lines is indicative of data that is to be written to a memory cell 12 coupled to the digit lines D, D*. An equilibration PMOS transistor 18 is also coupled between each pair of complementary digit lines D, D* to equalize the voltage between the digit lines D, D* prior to a memory read operation. Finally, a complementary PMOS bias transistor 20 is coupled to each digit line D, D* to lightly bias the digit lines D, D* to Vcc for reasons that will be explained. The current provided by each pair of bias transistors is controlled by a respective digit line load signal DLLN.\nA plurality of word lines WLI-WL4 are used to activate the memory cells 12 in the respective row of memory cells. The word lines WL1-WL4 are coupled to a respective inverter 30 each formed by a PMOS transistor 34 and an NMOS transistor 36 coupled in series between Vcc and ground. The gates of the transistors 34, 36 are coupled to each other and to a respective select line SEL WL1-SEL WL4.\nIn a read operation, the digit lines D, D* for all columns are equilibrated by driving an EQ* line low. An inverter 30 then drives a word line WL1-WL4 in a single row to an appropriate voltage, thereby coupling a memory cell 12 in each column to a respective pair of digit lines D, D*. The memory cell 12 in each column unbalances the digit lines D, D* to which it is coupled, and the respective sense amplifier 14 senses the polarity of the unbalance and provides an appropriate data signal.\nIn a write operation, a suitable voltage is first applied to a word line WL1-WL4 to couple the memory cells 12 in the respective row to a digit line D or a complimentary digit line D*. The write driver 18 for one or more columns then applies a differential voltage between the digit lines D, D* for respective columns, which is coupled to respective memory cells 12 for the activated row. The write driver 18 is then disabled for a xe2x80x9cwrite recoveryxe2x80x9d phase, and the word line WL1-WL4 is then deactivated so the memory cell 12 stores the polarity of the differential voltage. The bias transistors 20 are provided for the array 10 regardless of the type of SRAM cell used. However, in the event the memory cells 12 are loadless 4T memory cells, which are discussed further below, the current provided by the bias transistors 20 allow the memory cells 12 to continue to store the data, as also discussed further below.\nA typical memory cell shown in FIG. 2 is a conventional 6-transistor (6-T) SRAM cell 40. The SRAM cell 40 includes a pair of NMOS access transistors 42 and 44 that allow a differential voltage on the digit lines D, D*, to be read from and written to a storage circuit 50 of the SRAM cell 40. The storage circuit 50 includes NMOS pull-down transistors 52 and 56 that are coupled in a positive-feedback configuration with PMOS pull-up transistors 54 and 58, respectively. Nodes A and B are complementary inputs/output nodes of the storage circuit 50, and the respective complementary logic values at these nodes represent the state of the SRAM cell 40. For example, when the node A is at logic xe2x80x9c1xe2x80x9d and the node B is at logic xe2x80x9c0xe2x80x9d, then the SRAM cell 40 is storing a logic xe2x80x9c1xe2x80x9d. Conversely, when the node A is at logic xe2x80x9c0xe2x80x9d and the node B is at logic xe2x80x9c1xe2x80x9d, then the SRAM cell 40 is storing a logic xe2x80x9c0xe2x80x9d. Thus, the SRAM cell 40 is bistable, i.e., the SRAM cell 40 can have one of two stable states, logic xe2x80x9c1xe2x80x9d or logic xe2x80x9c0xe2x80x9d.\nIn operation during a read of the SRAM cell 40, a word-line WL, such as WL1-WL4 (FIG. 1), which is coupled to the gates of the access transistors 42 and 44, is driven to a voltage approximately equal to VCC to turn ON the transistors 42 and 44. The access transistor 42 then couples the node A to the digit line D, and the access transistor 44 couples the node B to the digit line D*. Assuming the SRAM cell 40 is storing a logic xe2x80x9c0xe2x80x9d, coupling the digit line D to the node A pulls down the voltage on the digit line D enough (for example, 100-500 millivolts) to cause the sense amplifier 14 (FIG. 1) coupled between the digit lines D, D* to read the SRAM cell 40 as storing a logic xe2x80x9c0xe2x80x9d.\nDuring a write operation of a logic xe2x80x9c1xe2x80x9d to the SRAM cell 40, for example, a logic xe2x80x9c1xe2x80x9d is applied to the digit lines D, D* as a differential voltage, and the word line WL is activated to turn ON the access transistors 42, 44. The transistor 42 then couples the logic xe2x80x9c1xe2x80x9d voltage of approximately VCC to the node A, and the transistor 44 couples the logic xe2x80x9c0xe2x80x9d voltage of approximately ground to the node B. The word line WL is finally deactivated to turn OFF the access transistors 42, 44, thereby allowing the SRAM cell 40 to continue storing the logic xe2x80x9c1xe2x80x9d.\nAlthough the 6-T cell 40 shown in FIG. 2 uses PMOS pull-up transistors 54, 58, it will be understood that other components (not shown), such as pull-up resistors (not shown), may be used in place of the pull-up transistors 54, 58.\nAnother typical SRAM cell is shown in FIG. 3. The SRAM cell shown in FIG. 3 is a conventional 4-transistor (4-T) loadless SRAM cell 60, where elements common to the SRAM cell 40 of FIG. 2 are referenced with like numerals or letters. The SRAM cell 60 is considered loadless because it uses a storage circuit 66 in which the loads formed by the pull-up transistors 54, 58 have been eliminated. Further, the NMOS access transistors 42 and 44 have been replaced with PMOS transistors 62 and 64, respectively. With the loadless 4-T SRAM cell 60 of FIG. 3, there are no pull-up transistors to maintain the drain of the OFF NMOS transistor 52, 56 at a voltage that is sufficient to turn ON the other NMOS transistor 52, 56. Instead, the access transistors 62, 64 are biased in their OFF states by conventional means with a voltage that causes leakage currents and/or subthreshold currents to be coupled from the digit lines D, D* through the access transistors 62, 64. These leakage currents and/or subthreshold currents maintain the voltage on the drain of the OFF NMOS transistor 52, 56, at a voltage that is sufficiently high to maintain the other NMOS transistor 52, 56 in an ON condition. In order to supply these leakage currents and/or subthreshold currents, the PMOS bias transistors 20 (FIG. 1) are controlled by the digit line load signals DLLN to supply currents to the digit lines D, D* when the memory cells 12 are not being accessed, as previously explained. However, the impedance of the transistors 20 must be sufficiently high that the digit lines D, D* in each pair can be driven low by the memory cells 12 and the write drivers 18.\nThe primary advantage of the 4-T SRAM cell 60 shown in FIG. 3 compared to the 6-T SRAM cell 40 shown in FIG. 2 is that the 4-T SRAM cell 60 uses only 4 transistors and is thus more compact. As a result, the 4-T SRAM cell 60 consumes less surface area on a semiconductor die.\nAlthough the loadless 4T SRAM cell 60 of FIG. 3 has the advantage of being more compact, it also has some disadvantages compared to the 6-T SRAM cell 40 of FIG. 2. These disadvantages primarily result from the need to supply the correct amount of leakage and/or subthreshold current through the access transistors 62, 64, and the need to use PMOS access transistors 62, 64 rather than NMOS access transistors 42, 44. Too little leakage and/or subthreshold current supplied to the storage circuit 66 may cause a data retention failure. If too much leakage and/or subthreshold current is supplied to the storage circuit 66, then the standby current limits of an array using the SRAM cell 60 may be exceeded.\nAnother problem resulting from the use of PMOS access transistors 62, 64 can be explained with reference also to FIG. 1. When power is initially applied to an integrated circuit containing the memory array 10, the digit lines D, D* can be driven to VCC before the word lines WL1-WL4 are driven high. With reference to FIG. 3, if the digit lines D, D* are at a high voltage when the voltage on the word line WL is low, the access transistors 62, 64 will be turned ON, thereby coupling the storage cell 66 to the digit lines D, D*. In fact, all of the SRAM cells 60 in the array 10 will generally be coupled to their respective digit lines D, D* under these circumstances. Although the leakage and/or subthreshold current drawn by any single SRAM cell 60 will be very small, the total current drawn by all of the SRAM cells 60 can be very large. For example, for a read current of as little as 100 microamperes (10xe2x88x924 amperes), the total current drawn by a 4 megabit SRAM array during power-up would be 400 amperes (10xe2x88x924* 4*106). Even though the current will not be this high in practice because of the finite current sourcing capability of the bias transistors 20, this amount of current is still far too much current to be drawn by SRAM memory devices.\nNote that the problem of excessive currents at power-up does not exist for the 6-T SRAM cell 40 shown in FIG. 2 because the NMOS access transistors 42, 44 will be OFF if the voltages of the word lines WL are less than the voltages on the digit lines D, D*. However, although not commonly in use, there may be circuit designs in which excessive power-up currents could be a problem even with NMOS access transistors 42, 44.\nThere is therefore a need for a system and method to limit the current drawn by SRAM arrays during power-up, particularly for arrays of SRAM cells having PMOS access transistors, such as loadless 4-T SRAM cells.\nAn array of SRAM cells arranged in rows and columns includes a wordline for each row of the array and a pair of complementary digit lines for each column of the array. Each of the SRAM cells has a pair of access transistors coupled to respective complementary digit lines for a respective column and a gate coupled to a wordline for a respective row. A bias circuit coupled to each of the digit lines is operable in either a normal mode or a power-up mode. In the normal mode, the bias circuit couples a bias current to the digit lines. In the power-up mode, the bias circuit maintains the access transistors non-conductive."}
{"text": "The drilling of oil wells has traditionally been dangerous and labor intensive, due in part to the potential for blowouts. A \"blowout\", which is an uncontrolled eruption of gas and oil from a well, is caused by a massive influx of formation fluid into the well bore at extremely high pressure. In many cases, blowouts occur at pressures sufficiently high to damage rig equipment and injure rig personnel. Blowout preventers (BOP's) are valves placed on top of a well to prevent blowouts. U.S. Pat. No. 4,535,852 to Boyadjieff and Krasnov describes the use of a blowout preventer valve in a top drive drilling system and is incorporated by reference in its entirety into the present disclosure.\nTop drive drilling systems (TDS's) rotate the upper end of a drill string directly by a drive system suspended from a traveling block, rather than using a traditional rotary table and kelly. U.S. Pat. No. 4,449,596 to Boyadjieff describes a TDS and is incorporated by reference in its entirety into the present disclosure. Two IBOP's (Inside Blowout Preventers) and a saver sub are typically threadedly connected to the drive shaft of a TDS. For example, often an upper IBOP can be threaded to the drive shaft, a lower IBOP can be threaded to the upper IBOP, and a saver sub can be threaded to the lower IBOP. The drill string is then threaded onto the saver sub. Thus, rotational motion is transferred from the drive shaft of the TDS, down through the upper IBOP, the lower IBOP and the saver sub, to the drill string.\nThe threaded connections between the drive shaft of a TDS, an upper IBOP, a lower IBOP, a saver sub and a drill string, must be torqued to a high level in order to survive the forces encountered in a conventional drilling operation, with the higher components being torqued significantly tighter than the lower components. Thus, these components are typically too tight to rely on the torque provided by a conventional TDS to break the connections. To provide the required high torques, U.S. Pat. No. 4,449,596 describes a special torque wrench, or \"pipe handler\", that provides additional force for making and breaking these tight connections. Such torque wrenches work well for \"breaking out\", or loosening, and \"making up\", or tightening, the connections between the upper IBOP, the lower IBOP, the saver sub and the drill string, but have not been able to be used for breaking out the upper IBOP. The problem is that the torque wrench has no way of reaching the main shaft of the TDS in order to keep it from rotating while applying torque to the upper IBOP.\nIn order to prevent blowouts, IBOP's are tested frequently to make sure they hold pressure. On average, a problem is found with an IBOP every two to three months, leading to its replacement. Traditionally, upper IBOP's have been removed manually using modified tong hangers which can weigh as much as 600 pounds. The hangers are taken off and the tongs are picked up with tugger lines. During this process, several strong people are required to hold the tongs level so they can bite into the upper IBOP and the main drive shaft of the TDS. Extreme lateral forces are exerted on the tongs, potentially causing them to break and injure the workers using them. This manual procedure takes an average of eight to ten hours of dangerous, difficult and frustrating labor. Clearly, it would greatly improve safety and efficiency if an automated method for replacing the upper IBOP could be provided."}
{"text": "Note that as a method for providing a metal auxiliary wiring for a transparent electrode of an electro-optical element, a method by which the metal auxiliary wiring is provided so as to overlap with an upper surface of the transparent electrode or a lower surface of the transparent electrode and to be electrically connected to the transparent electrode has been known (for example, see Patent Document 1).\nA structure in which an additional capacitor electrode provided for an active matrix substrate is formed of a conductive film having a light transmitting property of ITO, SnO2, or the like and an auxiliary wiring formed of a metal film is provided in contact with the additional capacitor electrode in order to reduce the electric resistance of the additional capacitor electrode has been known (see Patent Document 2).\nNote that it has been known that, as each of a gate electrode, a source electrode, and a drain electrode of a field effect transistor formed using an amorphous oxide semiconductor film, a transparent electrode of indium tin oxide (ITO), indium zinc oxide, ZnO, SnO2, or the like, a metal electrode of Al, Ag, Cr, Ni, Mo, Au, Ti, Ta, or the like, or a metal electrode of an alloy containing any of the above elements can be used; and, by stacking two or more of these layers, contact resistance may be reduced or interface intensity may be improved (for example, see Patent Document 3).\nNote that it has been known that, as a material of each of a source electrode, a drain electrode, a gate electrode and an auxiliary capacitor electrode of a transistor formed using an amorphous oxide semiconductor, a metal such as indium (In), aluminum (Al), gold (Au), or silver (Ag), or an oxide material such as indium oxide (In2O3), tin oxide (SnO2), zinc oxide (ZnO), cadmium oxide (CdO), cadmium indium oxide (CdIn2O4), cadmium tin oxide (Cd2SnO4), or zinc tin oxide (Zn2SnO4) can be used; and the same material or different materials may be used for the gate electrode, the source electrode, and the drain electrode (for example, see Patent Documents 4 and 5)."}
{"text": "The present invention relates to a tappet made of light metal such as Al alloy and a method of manufacturing the same.\nRecently, tappets used in a direct acting type valve operating mechanism in an internal combustion engine are made of Al alloy to lighten the valve operating mechanism and to increase allowable rotation speed of the engine. Al alloy tappets have lower strength and wear resistance than conventional steel tappets, so that wear resistant materials are mounted or coating layers are formed on the upper surface of the tappets which contact a cam and on the outer circumferential surface which contacts a cylinder head.\nFIGS. 8 and 9 show one example of a conventional tappet in which the upper surface which is slidably engaged with a cam is made to provide wear resistance. In FIG. 8, a groove la is formed on the upper surface of an Al alloy tappet body 1, and a wear resistant steel shim 2 is detachably engaged in the groove la. The shim 2 may be fixed by press-fit or casting.\nIn FIG. 9, a wear resistant cam receiving disc 3 is fixed on the upper surface of the tappet body 1 made of Al alloy by soldering.\nIn addition to the two examples, particles such as ceramics are dispersed on the whole surface of the tappet body to increase total strength and wear resistance.\nIn the tappet as shown in FIG. 8, effective surface area of the shim 2 which contacts the cam becomes smaller. Therefore, to obtain a desired lifting amount of. a valve, it is necessary to increase an outer diameter of the tappet body 1 and to mount a larger diameter shim 2. However, such structure increases the whole size of the tappet against lightening, and in addition, a cylinder head in which the tappet is mounted becomes larger to limit flexibility of planning in the engine. High accuracy mechanical processing must be made to the groove la, thereby increasing cost.\nIn the tappet in FIG. 9, the cam receiving disc 3 is fixed on the whole upper surface of the tappet body 1 to avoid decrease in effective surface area, but owing to connection of different materials, it is troublesome to solder them, thereby increasing cost.\nBoth the conventional tappets comprise two members, so that it is necessary to provide different size shims 2 and cam receiving discs 3 to make management of parts troublesome and not to attain lightening. Ceramic particles are dispersed on the whole tappet body, but it is difficult to disperse the particles uniformly, so that cost increases."}
{"text": "The invention relates to marine propulsion devices. More particularly, the invention relates to outboard motors, and still more particularly to outboard motor tilt tubes.\nA conventional outboard motor, such as the one disclosed in U.S. Pat. No. 4,710,141, comprises a transom or stern bracket mounted on the transom of a boat, and a tilt tube which is supported by the transom bracket and which extends along a generally horizontal tilt axis. Nuts threaded onto the opposite ends of the tilt tube prevent axial movement of the tilt tube relative to the transom bracket. A swivel bracket is mounted on the tilt tube for pivotal movement relative thereto about the tilt axis, and a propulsion unit is mounted on the swivel bracket for pivotal movement relative thereto about a generally vertical steering axis. A steering ram slidably housed within the tilt tube has one end connected to the steering arm of the propulsion unit via a power steering system, and has the opposite end connected to the core of a steering cable assembly. The sheath of the steering cable assembly is fixed to one end of the tilt tube by a nut threaded onto the end of the tilt tube.\nA conventional tilt tube extends a certain distance beyond one side of the transom bracket, which distance is only long enough to enable a retaining nut to be threaded onto the end of the tilt tube, but not long enough to enable the sheath of a remote control cable to be connected to that end of the tilt tube. The other end of the tilt tube extends a sufficient distance beyond the transom bracket to enable a retaining nut to be threaded onto that end of the tilt tube and to enable a remote control cable sheath to be secured to that end of the tilt tube by an additional nut.\nAttention is directed to the following U.S. patents:\n______________________________________ Nolan 3,207,117 September 21, 1965 Baba 4,735,165 April 5, 1988 Hickham 4,815,994 March 28, 1989 ______________________________________"}
{"text": "1. Field of the Invention\nThe present invention relates to a coupling structure for resin parts, and a method of manufacturing same. More particularly, the present invention pertains to such a coupling structure which may be used with or without tapping/tapping screws.\n2. Background Art\nAs a coupling structure for resin parts, there has been known a structure based on tapping in which the resin parts are fastened together using tapping screws as described, for example, in Japanese Patent Laid-Open No. 2012-61925 (see, particularly, paragraph [0028] and FIG. 6 of the document)."}
{"text": "Some sliding doors are designed, when closed, to recess completely into a frame that is installed within a wall between two opposing wall surfaces. Such doors are commonly referred to as “pocket” doors. Pocket doors provide the functionality of a hinge-mounted door with the additional advantage of space-savings. They have become increasingly popular, for example in small homes and in offices where square-footage is expensive.\nReferring to FIG. 1, a pocket door 5 is shown with an example of prior art hardware 20 that allows the user to manipulate the door, e.g. open, close, lock and unlock. Pocket doors 5 feature two broad surfaces 10 and an end surface 12. When the pocket door is in the open position, it is typically fully recessed in a frame contained in a wall (not shown). Pocket door hardware 20 can have a recessed area 30. The recessed area 30 allows easier manipulation of the door 5, e.g. as when opening or closing the door. Some versions of the hardware 20 feature a locking mechanism 35. Mounting screws 40 secure the door hardware 20 to the door 5.\nDoors with this type of hardware may in some cases be difficult to open from a fully closed, recessed position, particularly if the user's hands are full or the user has difficulty grasping objects."}
{"text": "Cable supports are used to organize and support medium voltage power distribution cables in underground manholes, vaults, and tunnels. Cable supports are also used to organize and support underground low voltage power cables and control cables, high voltage power transmission cables, and communication cables. Cable supports may also be used above ground and in areas other than underground manholes, vaults and tunnels.\nThese cables for electric power, control and communication lines are run underground in order to protect them from above-ground elements and from the interference and damage they would suffer when installed above the ground or on poles or structures. The underground environment may be less hostile in some ways, but the history of underground cables suggests that the underground environment is not benign. The environment in underground power and communications manholes is indeed harsh.\nWhile there may be fewer ultraviolet rays and less severe weather underground, and the temperature is more constant, moisture and humidity are always present. There are other considerations, such as the constant and higher danger from flooding, and underground pests that consider electrical insulation, and even steel, a tasty treat. Manholes may fill with water that is often contaminated with sewage, fertilizer runoff, tree roots, and chemicals, including caustic materials. Very harsh sea or salt water sometimes fills manholes. Many manholes are completely or partially filled with such contaminated water all of the time, except when pumped out for maintenance. Others fill periodically but are hot and have extremely high humidity, while still others fill and empty with ocean tides.\nAs noted, most power and communications manholes are partially or completely full of water some of the time or all of the time. The amount of water in a given manhole is influenced by location, surrounding conditions, drainage, and weather. Manholes located at higher grades generally will be filled with less water for a shorter period of time than those located at lower grades. Manholes located where the surrounding area has a high ground water level and/or a high amount of rain generally are filled with water to a higher level and more of the time than those located in areas that have a low surrounding ground water level and/or a low amount of rain. The water level in manholes located close to the ocean often changes with the tide, and the constantly-changing interface only increases the likelihood for corrosion. The condition of water in underground power and communications manholes occasionally is fresh and clean but most often is contaminated, as noted above, or is salt water, both of which can be very corrosive and also conductive.\nCommunication and power cables should be kept off surfaces, such as a floor or the ground, and should be organized and protected to the greatest extent possible. Cables are thus typically supported underground by racks that elevate cabling and keep the cabling off the ground, thus shielding the cables from at least some of the worst underground dangers. Racks for supporting cables must be able to withstand both heat and cold, all conceivable temperatures and humidities in every combination. In addition, the racks must be able to support very heavy loads from power and communication cables. The racks themselves are preferably supported, e.g., attached to a wall, rather than free-standing structures. Thus, the racks will have penetrations, or stress concentrators, to deal with, in these hot, humid, and stressful environments, along with the high loads expected from supporting cabling. The walls themselves may have penetrations for supporting bolts, pins or other fasteners used to secure the racks in place. The walls, such as concrete walls or other structures, will also be in intimate contact with the racks, adding their chemical potential for corrosion to the racks.\nAll these stresses combine to make the underground a challenging environment for cable racks. For the most part, existing cable supports used in underground manholes, vaults, and tunnels are manufactured using steel stampings, steel forms, or steel weldments. They may also be ductile iron castings. After the supports are stamped, formed, welded, or cast, they are hot dip galvanized in an effort to prevent corrosive deterioration. The steel arms and posts are bonded together and grounded in an attempt to prevent corrosion. Eventually, the galvanized coating is consumed and the steel racks may oxidize or corrode away, leaving the power and communications cables without support.\nTwo phenomena, galvanic corrosion and stray current corrosion, occur in flooded underground manholes to cause this deterioration. Galvanized steel cable supports are very vulnerable to both galvanic and stray current corrosion and often become severely corroded to a point that they will no longer support the cables in a very short period of time.\nGalvanic corrosion is an electrochemical process in which one metal, the anode, corrodes preferentially when in electrical contact with a different type of metal, the cathode, and both metals are immersed in an electrolyte. In flooded underground power and communications manholes the galvanized steel cable supports are the anodic sites of the galvanic corrosion reaction. Cathodic parts in the manhole, parts made from more noble metals such as stainless steel, may be damaged in the galvanic corrosion process due to generation of electrolytic hydrogen on their surfaces causing hydrogen embrittlement. Stray current corrosion of underground power and communication cable supports is usually caused by power and communications manholes being located in the vicinity of electric rail tracks, pipe lines that are cathodicly protected or the like.\nUnderground galvanized steel cable supports that are severely corroded and can no longer support the cables result in power and communications interruptions and a safety hazard to technicians who enter the manhole. Another safety issue is that galvanized steel cable supports are conductive. If a power cable's insulation is compromised and the electrified conductor contacts a galvanized steel cable support, the cable support is energized. If a technician inadvertently touches the energized cable support he may be electrocuted.\nWhat is needed are safer cable racks better able to withstand the environment and better able to tailor themselves to a greater variety of situations, for fewer stresses, and for longer service."}
{"text": "In some instances it is beneficial to utilize an apparatus comprising a compensating interferometer that modulates the phase of an input light signal to produce an output signal that is a superposition of an input signal and a phase shifted version of the input signal. One such instance is when a compensating interferometer is used to manipulate the output of a sensor array as described in U.S. Pat. No. 5,987,197. However, known compensating interferometers lack functionality that can be advantageous in processing optical signals."}
{"text": "Hepatitis C virus (HCV) has plus single-stranded RNA genome. The gigantic precursor protein of about 3,000 amino acids translated from the HCV RNA is cut into a virus particle-forming core protein and envelope protein, and other non-structural proteins by host cell-and virus-derived protein degrading enzymes.\nThe core protein is known to travel into the host cell nucleus, in addition to forming the virus particle. Recently, a strong involvement of the core protein in the onset of liver cancer in which the core protein modulates host cell functions in a variety of ways has been reported. Non-Patent Literature 1 describes, for example, HCV core protein-expressing transgenic mice developing liver cancer via fatty liver with high probability. In an effort to elucidate the molecular mechanism underlying the onset of liver cancer caused by the core protein, it has been common to identify host proteins that interact with the core protein, and analyze the functions of such proteins.\nMany host proteins that interact with the core protein have been reported, including p53 (Non-Patent Literature 2), RNA helicase (Non-Patent Literature 3), STAT3 (Non-Patent Literature 4), and PA28γ (Non-Patent Literature 5).\nNon-Patent Literature 5 reports that the core protein interacts with PA28γ, and localizes in the nucleus. Of note, PA28γ is known as a proteasome regulatory protein localized within the cell nucleus that interacts with 20S proteasome and improves the peptidase activity thereof. This literature also reports that the 44 to 71 amino acid region of the core protein is involved in both PA28γ binding and localization in the nucleus, and that the core protein is subject to PA28γ-dependent degradation. However, the specific mechanism underlying the involvement of the core protein in the onset of liver cancer is not fully elucidated in any of the foregoing literatures.\nPatent Literature 1 discloses a screening method of a drug useful for the prevention and/or treatment of hepatitis C virus-associated disease, the method including the step of evaluating the inhibitory activity on the protein interaction between the hepatitis C virus core protein and PA28γ, and a preventive and/or therapeutic drug for hepatitis C virus-associated disease obtained by the method. However, this literature is concerned with suppressing the PA28γ-induced onset of a disease arising from the hepatitis C virus-derived core protein, specifically suppressing the progression of hepatitis conditions to fatty liver, cirrhosis, and liver cancer, and does not clarify functions that prevent hepatitis C virus infection, or suppress HCV growth. Thus, unlike antiviral drugs, the hepatitis C virus-associated disease preventive and/or therapeutic drug of this literature is not intended to suppress or reduce virus amounts, and only alleviates symptoms. Therefore, continuous administration is used, and clinical development is expected to take a long time because the persistent infection takes 10 to 15 years before developing into hepatitis C virus-associated diseases such as cirrhosis and liver cancer.\nIt would, however, be possible to cure HCV infections if the HCV could be cleared with the use of antiviral drugs.\nConsidering that the current mainstream antiviral drugs target virus replication, it is considered possible to further improve anti-viral effects with an antiviral drug that uses a new mechanism to prevent viral infection or suppress virus growth."}
{"text": "Breast cancer is a leading cause of death in women. While the pathogenesis of breast cancer is unclear, transformation of normal breast epithelium to a malignant phenotype may be the result of genetic factors, especially in women under 30. Miki, et al., Science, 266: 66-71 (1994). However, it is likely that other, non-genetic factors also have a significant effect on the etiology of the disease. Regardless of its origin, breast cancer morbidity increases significantly if a lesion is not detected early in its progression. Thus, considerable effort has focused on the elucidation of early cellular events surrounding transformation in breast tissue. Such effort has led to the identification of several potential breast cancer markers. For example, alleles of the BRCA1 and BRCA2 genes have been linked to hereditary and early-onset breast cancer. Wooster, et al., Science, 265: 2088-2090 (1994). The wild-type BRCA1 allele encodes a tumor supressor protein. Deletions and/or other alterations in that allele have been linked to transformation of breast epithelium. Accordingly, detection of mutated BRCA1 alleles or their gene products has been proposed as a means for detecting breast, as well as ovarian, cancers. Miki, et al., supra. However, BRCA1 is limited as a cancer marker because BRCA1 mutations fail to account for the majority of breast cancers. Ford, et al., British J. Cancer, 72: 805-812 (1995). Similarly, the BRCA2 gene, which has been linked to forms of hereditary breast cancer, accounts for only a small portion of total breast cancer cases. Ford, et al., supra.\nSeveral other genes have been linked to breast cancer and may serve as markers for the disease, either directly or via their gene products. Such potential markers include the TP53 gene and its gene product, the p53 tumor supressor protein. Malkin, et al., Science, 250: 1233-1238 (1990). The loss of heterozygosity in genes such as the ataxia telangiectasia gene has also been linked to a high risk of developing breast cancer. Swift, et al., N. Engl. J. Med., 325: 1831-1836 (1991). A problem associated with many of the markers proposed to date is that the oncogenic phenotype is often the result of a gene deletion, thus requiring detection of the absence of the wild-type form as a predictor of transformation.\nOf interest to the present invention are reports that the protein content of the nuclear matrix in breast epithelia may provide a marker of cellular growth and gene expression in those cells. Khanuja, et al., Cancer Res., 53: 3394-3398 (1993). The nuclear matrix forms the superstructure of the cell nucleus and comprises multiple protein components that are not fully characterized. The nuclear matrix also provides the structural and functional organization of DNA. For example, the nuclear matrix allows DNA to form loop domains. Portions of DNA in such loop domains have been identified as regions comprising actively-transcribing genes. Ciejek, et al., Nature, 306: 607-609 (1982). Moreover, the organization of the nuclear matrix appears to be tissue-specific and has been associated with so-called transformation proteins in cancer cells. Getzenberg, et al., Cancer Res., 51: 6514-6520 (1991); Stuurman, et al., J. Biol. Chem., 265: 5460-5465 (1990).\nProteins and steroid hormones thought to be involved in transformation are associated with the nuclear matrix in certain cancer cells. Getzenberg, et al., Endocrinol. Rev., 11: 399-417 (1990). It has been suggested that changes in the composition or organization of nuclear matrix proteins may be useful as markers of growth and gene expression in breast tissue. Khanuja, et al., Cancer Res., 53: 3394-3398 (1994). However, Khanuja did not identify any specific proteins for use as cancer markers.\nThere is, therefore, a need in the art for specific, reliable markers that are differentially expressed in normal and transformed breast tissue and that may be useful in the diagnosis of breast cancer or in the prediction of its onset. Such markers and methods for their use are provided herein."}
{"text": "The density of a substance is expressed as a ratio of mass m to volume V, or m/V This is a physical property of a material which relates to composition, level of impurities, and mixtures, and can be an indicator of hidden features such as voids. In the case of compressible media, such as closed-pore solids, the bulk density is a function of hydrostatic pressure, since the volume changes but the mass remains constant.\nThere exist two popular methods for determining the density of a solid: (1) by comparison of the sample density with the densities of substances of known value, usually by hydrostatic weighing in two different fluids of known and substantially different density (Archimedes principle), and (2) by the independent measurement of mass and volume of the sample.\nConsidering the first method, the weight of an object is measured in two liquids having significantly different densities. The measured weight, also referred to herein as apparent weight, is reduced from the true weight due to buoyancy forces acting on the object Thus, the apparent weight is the true weight minus the buoyancy force, where the buoyancy force is equal to the weight of the liquid displaced by the object. When in the higher density liquid, the buoyancy force acting on the object is greater, and the apparent weight of the object is less. The sensitivity of the method ultimately relies on the range in available liquid densities. In one method of using hydrostatic weighing techniques, one measures the apparent weight of an object in alcohol and then in water, where the alcohol and the water have densities of 0.791 and 1.0 g/cc, respectively. It is also common practice to weigh the object in air and then in water or alcohol, thus using air as the first medium. In that case it is common practice to assume the weight in air to be the “true weight” of the object. The specific reasons for using a liquid as one of the mediums in this technique are (i) to obtain a large difference in density between the two fluids, and (ii) to increase the effect of buoyancy forces.\nFor very accurate measurement of density, there are several experimental problems which are typically ascribed to the hydrostatic weighing method using two liquids. First, the method suffers from the necessity of weighing an object in liquid. Strictly as a practical matter, this requires suspending the object via a tether or thin wire in the liquids. Second, related to the first, is the fact that surface tension forces affect the measurement as the liquid meniscus either pulls the tether down into the liquid or pushes it up into the adjacent gas (typically air), depending on whether the liquid wets the tether easily. Third, the density of the liquid is affected by dissolved gases in the liquid. Since the effect of trapped gas is to change the actual density of the liquid, efforts must be made to eliminate the trapped gas. Fourth, results will vary due to bubbles of trapped gas on irregular sample surfaces of the object. The bubbles that cling due to surface tension displace liquid and affect the measured buoyancy forces.\nThe second common approach to determine density requires independent determination of both the mass and the volume. One measures the mass of the body using conventional state-of-the-art balances common to most laboratories. Commercial devices exist for performing this step to very high precision and accuracy. The volume is determined independently. If the sample is of a uniform geometry, it may be possible to calculate the specimen volume based on measurable dimensions. In the more general case where samples are of irregular shape, the volume is determined by a method commonly referred to as pycnometry. For reasonably sized samples on the order of 0.5 cubic centimeters and larger, commercial pycnometers are available for determining volume to 0.02%. Pycnometers typically consist of two chambers connected by means of a pathway for a gas to move and a valve which can isolate the two chambers. The exact volume of one of the chambers must be known apriori. The second chamber is of arbitrary, but similar size. The first chamber, of known volume, is pressurized using a gas such as helium to a predetermined pressure. The second chamber is initially empty and is evacuated by means of a vacuum pump. By means of valves, the two chambers are then isolated from the gas source and from the vacuum pump leaving the first chamber at an elevated pressure with helium and the second chamber under vacuum. The valve in the passageway connecting the two chambers is then opened and the pressurized gas is allowed to expand from the first chamber into the second chamber, and the pressure achieves a new equilibrium value by virtue of the increased volume occupied by the gas. It is a straightforward calculation to determine the volume of the unknown chamber using the initial helium pressure in the first chamber, the volume of the first chamber, and the final pressure. The steps are then repeated with the sample of interest being placed into the second chamber. The newly calculated volume of the second chamber represents the remaining volume of the second chamber not occupied by the sample.\nAlthough the field of pycnometry is well established for accurately determining the volumes of solids of reasonable sizes, the state-of-the-art is limited by several factors in attempts to extrapolate to smaller samples. There are many industries where large samples are not always available. Some specific applications would include high-temperature superconducting wires, samples pertaining to the study of irradiation, and porous membranes used for delivering and mixing gases, such as in the fuel cell applications. The volume of such samples is often much smaller than that required by pycnometers. For accurate measurements, the sample should occupy a significant fraction of the chamber volume, e.g. 50-60% of a 1 cubic centimeter chamber.\nThere are several other limitations to the pycnometry method for determining density. First, if safeguards are not included, temperature variations of the gas due to room temperature fluctuations may affect the pressure, and hence the density, of the pressurized gas. Second, the gas-comparison pycnometer described above does not work if there are any leaks in the system. The ability of the technique to work depends strongly upon the number of gas atoms remaining constant before and after the gas expands into the second chamber. Third, the chamber door, when closed, must close in a repeatable fashion such that the volume of the chamber is exactly the same every time the door is opened and closed. Fourth, the mechanism requires a vacuum pump. Fifth, the true volume of the pressurized chamber must be known to better tolerances than the desired accuracy of sample volume. Finally, the mass must be measured by a separate device.\nThus, there are significant limitations associated with the prior art The current invention offers significant improvements over the prior art, as will become apparent in the following description of invention."}
{"text": "Rapid advances in electronics and communication technologies, driven by immense customer demand, have resulted in the widespread adoption of electronic transactions, record keeping, and data sharing. As one example, e-currencies, such as Bitcoin, have displaced paper currencies in millions of transactions per year. Improvements in tracking and/or knowledge management attached to such electronic transactions will continue to increase the features and options available to operators engaging in electronic transactions."}
{"text": "Wind turbine blades are a means of converting power from wind into mechanical power to propel a generator of a wind turbine. Wind turbine blade design is mainly conducted under consideration of aerodynamic and mechanical considerations, i.e. the airfoil shape or profile is designed as a compromise between aerodynamic optimization and mechanical characteristics, for example strength, for various wind conditions. An optimal blade of a wind turbine exhibits a low cut in wind speed, good and efficient aerodynamic performance for various wind speeds, which means for low to high wind speed, minimal loads on blade and wind turbine structure, i.e. under turbulent and high wind speed conditions, and low acoustic noise.\nIt is for instance important that the blade can resist high forces and stresses close to the hub. Blades are therefore thick and wide in the region of a root section close to the hub. At the root, the blade is usually narrow and tubular to fit onto the hub and to provide sufficient strength. The blade profile becomes thinner and thinner as to obtain acceptable aerodynamic properties. The lift force will increase as the speed increases towards a tip of the blade, caused by a larger diameter at the same rotational frequency toward the tip. Decreasing the chord width towards the tip will contribute to counteract this effect. The blade tapers from a point somewhere near the root towards the tip. Furthermore the blade is twisted along its axis to account for a change in direction of the airflow to the wind resulting from rotation. The speed of a blade section is increasing the further it is located towards a tip of the blade.\nNevertheless, conventional blades of wind turbines have a static airfoil and thus exhibit limited possibilities to adjust for wind conditions. The only possible means of adjustment and optimization of the aerodynamic properties for a wind speed and turbulence is the adjustment of a pitch angle for the whole blade. From an aerodynamic point of view an optimal blade of a wind turbine would not only comprise a variable pitch angle for each section of the blade but also comprise of an adjustable airfoil to account for different wind and turbulence conditions.\nIt is commonly known to enhance and optimize the performance of wind turbine blades with devices added onto the wind turbine blade. Such devices are among others active or passive components such as flaps, vortex generators or stall strips. The actuation of flaps can for instance be conducted with electricity, hydraulic or piezoelectric means. The bending of a blade is considered to stay in a relation to blade loading and wind conditions. A passive solution is known, where a shape and/or angle of a flap relative to a blades chord line changes depending on the bending of a wind turbine blade.\nAnother difficulty with wind turbine blades is the occurrence of acoustic noise when in operation. The tip speed of a blade in operation is for instance 80 m/s. One means to reduce the noise is the attachment of a serrated plate in the region of the trailing edge projecting over the trailing edge. The EP 1314885 discloses valuable information on wind turbine blades technology and performance.\nGenerally, it is desirable to improve the aerodynamic performance and efficiency of wind turbine blades, i.e. in low wind speed conditions. Moreover, the mechanical loading of the blade may be to be minimized. A secondary factor is the reduction of acoustic noise of the wind turbines blades when in operation.\nFurther relevant state of the art is disclosed in EP 1623111 B1, US 2011/0116927 A1, WO 2004/088130 A1 and EP 2034178 A2. The EP 1314885 discloses an apparatus improving the efficiency of a wind turbine with a panel connected to the trailing edge of the wind turbine blade. EP 2034178 A2 discloses a fairing plate to avoid an air gap when the flap is deflected, but which has no mechanical function."}
{"text": "Gas turbine engines often include so-called \"start injectors\" that are used when initiating operation of the turbine. Start injectors operate on the principal of pressure atomization and as a consequence, in relatively small turbine engines in airborne environments whereat fuel flows at high altitudes during starting may be quite low, high pressures are required to achieve the pressure atomization required of a start injector. Air blast atomization is not a viable alternative during starts when the turbine is rotating at a minor percent of its rated speed. At the same time, at such low speeds, it is difficult with available fuel pumps to generate the necessary fuel pressure.\nTo meet these difficulties, conventional start injectors have extremely small orifices to provide the desired atomization making them precision formed parts. They are thus costly to manufacture and at the same time, because of the very small orifices employed, they are prone to plugging.\nIn the above-identified application of Shekleton, there is disclosed a unique fuel injector construction for a turbine engine that is of low cost construction and which may advantageously be employed both as a start injector and as a main fuel injector and which does not require high fuel pressures to achieve good atomization. The present invention is directed to improving upon the invention disclosed in the previously identified Shekleton application."}
{"text": "FIG. 1 is an overview of an open water workover system 1 according to prior art. The system 1 is a set of equipment which is installed on top of a Christmas tree (XT) 2 and allows safe intervention operations in a hydrocarbon containing well (not shown). A well control package (WCP) 3 also named lower riser package (LRP) is connected to the top of Christmas tree 2. At the top of the lower riser package (LRP) 3, an emergency disconnect package (EDP) 4 is attached. A workover riser 5 is extending from a drilling rig 6 on the sea level 7 down to the emergency disconnect package 4.\nFIG. 2 shows a typical setup for a prior art horizontal Christmas tree 10 with a monobore workover system 11. In this setup a monobore or workover riser 12 extends from the rig (only the lower part of the riser is shown) down to a monobore emergency disconnect package (EDP) 13 and lower riser package (LRP) 14. The lower riser package 14 is connected to the horizontal Christmas tree 10 through a Christmas tree adaptor connector (XTAC) 15. The Christmas tree 10 is in the other end connected to a wellhead 16 and a tubing hanger 17 sits in the Christmas tree in conventional manner for horizontal Christmas trees 10. The emergency disconnect package 13, the lower riser package 14 and the Christmas tree adaptor connector 15 are constituting the monobore workover system 11. The system is referred to as monobore due to a single, straight main bore 18a that runs down the centre of the system and is adapted to correspond with a main bore 18b in the horizontal Christmas tree 10. An annulus bore 26 encloses the main bore 18b in the horizontal Christmas tree 10. The horizontal Christmas tree 10 has a horizontal production outlet 19 where the production flow is routed. This outlet 19 is closed during the workover process by valves 20. When the workover system 11 is attached to the Christmas tree 10, tools (not shown) can be lowered and hoisted through the riser 12 and main bore 18a, 18b in order to perform different tasks in the well.\nIn FIG. 2 an annulus line 21 is arranged beside the main bore 18a in the workover system 11. The annulus line 21 starts as a hose (not shown) from topside and enters the emergency disconnect package at the top part 22 before taking a varied path down through the system and into the Christmas tree 10. There are several valves in the annulus line 21 and several crossover points 23 from the annulus line 21 into the main bore 18a, creating circulation paths that can be used for various operations. The line 21, valves 24 and crossover points 23 can together with valves 25 in the main bore 18a for instance be used for flushing of hydrocarbons from the riser 12.\nFIG. 3 shows a typical setup for a prior art vertical Christmas tree 31 with a similar monobore workover system 11 as in FIG. 2, but includes a passive bore selector 32. The vertical Christmas tree has a configuration with a production bore 33 and annulus bore 34. When using a monobore workover system on the vertical Christmas tree 31 it is therefore required to use a bore selector technology, either a passive or an active bore selector, to match the configuration of the vertical Christmas tree. The workover system 30 in FIG. 3 has a similar structure than the monobore workover system 11 in FIG. 2, except that there is arranged a passive bore selector 32 between the lower riser package and the vertical Christmas tree 31. The main bore 18a in the lower riser package 14 is adapted to correspond with the production bore 33 in the vertical Christmas tree 31 through the passive bore selector 32. The same applies to the annulus line 21 in the lower riser package 14 which is adapted to correspond to the annulus bore 34 in the vertical Christmas tree 31. The Christmas tree 31 is in the lower end connected to a tubing hanger 35 and a wellhead 36 in a conventional manner for vertical Christmas trees 31\nWith use of the passive bore selector 32 the whole riser system must be pulled to the surface in order to change bore selector.\nThe present invention is applicable both on Vertical Christmas Trees (VXT) as well as other types of Christmas trees. In the Figures illustrating the invention there are shown a vertical Christmas tree and the invention will be further described for this use, but the invention is not limited only to vertical Christmas tree, but is also applicable for other types of Christmas trees, such as hybrid Christmas trees, horizontal Christmas trees or other possible types of Christmas trees.\nTo retrieve or deploy a Vertical XT, it is required to set plugs in the Tubing Hanger—one plug in the main bore and one plug in the annulus bore.\nTraditionally this has been done using either a dual bore workover system or a monobore workover system with active or passive bore selectors, as described in FIG. 3.\nA dual bore work over system has several limitations—especially on the Emergency Disconnect Package (EDP) and Lower Riser Package (LRP) design. The two bores need to be straight to allow straight access to the Christmas tree and the tubing Hanger. Different bore configurations in the XT will require the same corresponding bore configuration in the EDP/LRP, hence making standardization of the EDP/LRP design difficult. A dual bore system also requires an extra 2″ riser in addition to the normal main bore riser—something that can lead to a reduced operational window along with increased wellhead fatigue. Some special joints, like a safety joint or a heave eliminator are very complex to make for a dual bore system.\nUsing a monobore workover system requires the use of bore selector technology. A passive bore selector is an alternative connector assembly that can be used that offsets the whole lower workover riser package (WRP) stack, i.e. EDP and LRP to align the main bore over the annulus bore to give the required access. A passive bore selector has traditionally been chosen due to the reliability and the small size, but the operational efficiency is greatly reduced since the whole riser stack must be pulled to the surface in order to change the bore to be intervened.\nAn active bore selector allows wireline tools to be re-directed from the main bore in the EDP to either the main or the annulus bores in the LRP, with a path that is straight enough for the wireline tools to navigate. Using an active bore selector has shown to add complexity to an already complex system, and it is required to use a flexible jarring tool for setting of the annulus plug due to the “dog leg” going from main to annulus bore. The risk for getting stuck with the jarring tool, or failing with the jarring due to lack of force transfer is greatly increased through a “dog leg”. In addition, an active bore selector will typically add weight and height to a system that is already above the weight/height limitations stated in commonly used standards.\nWith both of the above arrangements, this creates a heavier, more complicated and more expensive system design. Therefore, the invention proposed is aimed to solve this problem by creating a separate tool that can be deployed to perform the annulus plugging/removal operation when necessary."}
{"text": "1. Field of the Invention\nThe present invention relates to image-forming devices including an independent control panel which is separate from a main body of the image-forming device. In particular, the present invention relates to a copying device including an independent control panel having a display function, arranged so as to minimize the size of the copying device, and particularly to minimize the floor space required for installation of the copying device, including the main body itself and an operational area of the copying device.\n2. Description of the Related Art\nA copying device, as a known image-forming device, is shown in FIGS. 7 and 8. Other known image-forming devices include printers, facsimiles, combinations of printers and facsimiles, and the like.\nA copying device 100 shown in FIGS. 7 and 8 is provided with an automatic document-feeder 102 at an upper part of a main body 101 of the copying device 100. The automatic document-feeder 102 automatically feeds original sheets to an image reader 103, as shown in FIG. 8. The automatic document-feeder 102 includes a mounting part 102a for setting the original sheets thereon. The image reader 103 optically reads the original sheets or books that are placed upon the image reader 103. The image reader 103 is disposed at the uppermost part of the main body 101, which opposes a lower face of the automatic document-feeder 102.\nThe automatic document-feeder 102 can pivot on hinges disposed at the rear side thereof. This allows the automatic document-feeder 102 to be maintained open at any desired angle with respect to the main body 101, so that an original sheet or a book can be set on the image reader 103 for reading.\nAn original-sheet discharge-tray 102b for stacking original sheets which have been read and discharged by the automatic document-feeder 102 is provided at the right side of the upper part of the main body 101. The main body 101 is provided with a display 106a at the front side of the automatic document-feeder 102 and at the upper part of the main body 101.\nA sheet-supply tray 105 for stacking sheets for printing, when the sheets are manually supplied to the main body 101 by an operator, protrudes from the left side of the main body 101.\nWhen the operator operates the copying device 100 by using a control panel 106, the image reader 103 reads data from an original sheet or a book, an imaging unit (not shown) provided in the main body 101 transfers an image data corresponding to the read data to a sheet which was supplied from a sheet cassette 107 or the sheet-supply tray 105, and the sheet is discharged to a sheet-discharge tray 110.\nThe control panel 106 of the known copying device 100 is positioned at the front side of the upper part of the main body 101, that is, between the automatic document-feeder 102 and the operator. The control panel 106 preferably is positioned toward the front side of the copying device 100, because if it is disposed toward the rear side of the copying device 100, the control panel 106 will interfere with the easy opening and closing of the automatic document-feeder 102 or setting of the original sheets or books on the image reader 103.\nTherefore, the control panel 106 conventionally has been positioned toward the front side. Since the space provided to install the control panel 106 is limited at the front side, however, the size of the control panel 106 has had to be reduced.\nIn recent years, however, the size of the control panel 106 of known copying devices has increased. Currently, a large liquid crystal monitor including a touch panel often is provided in order to cope with the trend toward multiple functionality of the copying device and improved maneuverability thereof. As a result, an enlarged control panel 106 displaces the automatic document-feeder 102 and the image reader 103 toward the rear side of the main body 101, whereby the original sheets or books are not easily set on the image reader 103 and the size of the copying device 100 is increased.\nIn particular, it may be particularly difficult for an operator in a wheelchair, who operates the copying device 100 from a position lower than that of an operator who operates the same while standing, to set original sheets or books on the image reader 103 and handle the automatic document-feeder 102.\nMoreover, it is difficult for an operator in a wheelchair to view a control panel 106, which is substantially horizontal to the upper face of the main body 101 (the control panel 106 has poor visibility for the operator in a wheelchair).\nThe visibility of the control panel 106 may be improved by providing an operation surface of the control panel 106 that is substantially vertical, and by providing the control panel 106 toward the front side of the main body 101 of the copying device 100.\nHowever, when the control panel 106 is disposed toward the front side of the main body 101, the control panel 106 and a supporting member for the control panel 106 may interfere with the opening and closing of the automatic document-feeder 102 and with the opening and closing of a cover 104 of the main body 101 of the copying device 100. However, if the control panel 106 is disposed toward the rear side of the main body 101 of the copying device 100, remote from the operator, operability and maneuverability of the copying device 100 is reduced.\nCopying devices that attempt to overcome the above drawbacks are disclosed in, for example, Japanese Patent Laid-Open Nos. 9-6185 and 11-119498. The copying device disclosed in Japanese Patent Laid-Open No. 9-6185 is provided with a control panel 121 having a display 121a, as shown in FIG. 9, which is movable to a position appropriate for operation. The copying device disclosed in Japanese Patent Laid-Open No. 11-119498 is provided with an inclined control panel 131, as shown in FIG. 10, of which the angle of inclination is variable.\nEach of copying devices 120 and 130 shown in FIGS. 9 and 10, respectively, is provided with sheet cassettes 107 which can be drawn from the main body 101. The copying device 120 is provided with a cover 104 for opening and closing the main body 101. The automatic document-feeder 102 is provided with a mounting part 102a for setting original sheets thereon.\nThe copying devices 120 and 130 are provided with control panels 121 and 131, respectively, and displays 121a and 131a, respectively, which are positioned at the respective front sides of the automatic document feeders 102. Since the control panel 121 or 131 is positioned at the front side, the control panel 121 or 131 interferes with the opening and closing of the automatic document-feeder 102. Therefore, the display 121a of the control panel 121 must be moved or the angle of inclination of the control panel 131 must be changed each time the automatic document-feeder 102 is opened or closed.\nMoreover, it is difficult to position the automatic document-feeder 102 further toward the front side, and it is necessary to form a space around a main body 122 or 132 of the copying device 120 or 130, respectively, in which the display 121a or the control panel 131, respectively, can move. That is, there is a drawback in that the space around the copying device 120 or 130 is inefficiently used, because, when the device is operated, a larger operational space is required in addition to the space which is occupied by the copying device 120 or 130 when it is operated.\nAccordingly, it is an object of the present invention to provide an image-forming device including an independent control panel provided with a display function, in which an installation space of the device, including an operational space of a main body of the image-forming device, is minimized.\nIt is another object of the present invention to provide an image-forming device including an independent control panel, in which the control panel is easily accessible by an operator, and in particular, an operator confined to a wheelchair, without requiring excessive reach to the control panel, when the operator operates the control panel at the front side of the image-forming device.\nTo these ends, according to an aspect of the present invention, an image-forming device comprises a main body; a protrusion extending from a side surface of the main body; a support member provided either on the protrusion or on the protrusion and the main body; and a control panel including a display. The control panel is separate from the main body and supported by the support member toward a front side of the protrusion.\nThe support member of the image-forming device according to the present invention may extend above the protrusion from a rear side of the protrusion toward the front side thereof and support the control panel at the extended distal end of the support member.\nThe control panel of the image-forming device according to the present invention may be supported by the supporting member such that the control panel is capable of pivoting, rotating or turning in at least one of a vertical direction and a horizontal direction, by any arbitrary angle.\nThe protrusion of the image-forming device according to the present invention may include a control member which provides control in connection with at least one operation/function of the control panel, such as copy start, copy count, and the like.\nThe control member of the image-forming device according to the present invention may be provided at the front side of the protrusion.\nThe protrusion of the image-forming device according to the present invention may comprise a tray for stacking original sheets or sheets for printing images to be supplied to or discharged from the main body.\nThe protrusion of the image-forming device according to the present invention may comprise a mounting table for receiving thereon or therein operation manuals.\nThe control panel of the image-forming device according to the present invention is separate from the main body and supported by the support member toward the front side of the protrusion, where the support member is provided either on the protrusion or on the protrusion and the main body, whereby the maneuverability, accessibility and visibility of the control panel is improved.\nThe control panel of the image-forming device according to the present invention is supported and is disposed using a protrusion which is necessary for operation of the image-forming device. Therefore, it is not necessary to provide an additional exclusive mounting space, whereby the size of the image-forming device can be reduced.\nSince the control panel does not protrude significantly from the front face of the main body, the image-forming device can be reduced in size and the space necessary for operation can be reduced to a minimum, whereby the image-forming device can be installed in a more compact space.\nFurther objects, features and advantages of the present invention will become apparent from the following description of the preferred embodiments with reference to the attached drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to pumps which operate on the gerotor fluid displacement principle and more particularly to improved pump constructions for implementing this principle by means of magnetic coupling into a sealed pumping compartment.\n2. Description of Prior Art\nIn order to isolate the handled fluid from the environment of the pump motor, a variety of pump configurations have been devised wherein there is a magnetic field coupling between a motor shaft and a drive shaft for the driven pump member. U.S. Pat. Nos. 2,970,548 and 2,996,994 are exemplary of concentric and axial approaches for such magnetic couplings. U.S. Pat. No. 4,526,518 discloses an application of the axial coupling approach to supply drive to the shaft of a male element of a gerotor pump.\nAs is described in the aforementioned '518 patent, gerotor pumps offer advantages as to pumping efficiency and low friction. Other advantages of such pumps include quietness in operation and self-priming modes of operation. While the '518 pump does successfully implement the magnetic drive of gerotor elements, it is relatively complex in construction and has a fairly large number of parts."}
{"text": "The present invention relates generally to a content recording/reproducing apparatus and method for recording and reproducing content, a storage medium, and a computer program. More particularly, the present invention relates to a content recording/reproducing apparatus, a storage medium, and a computer program which allow general viewers (listeners) to record and reproduce content distributed by broadcasting or other means.\nTo be more specific, the present invention relates to a content recording/reproducing apparatus and method, a storage medium, and a computer program which record and reproduce content limited in recording retention period or reproducing period for the purpose of copyright protection for example. More particularly, the present invention relates to a content recording/reproducing apparatus, a storage medium, and a computer program which allow users who view content limited in recording retention period or reproduction period to easily understand a mechanism of the time limits of the recording retention period and the reproduction period.\nToday, the advances in digital technology allow the storage of AV data composed of audio and video and in immense quantities without deterioration in quality. For example, hard disk drives capable of storing several tens of gigabytes or more are now available at comparatively low prices and recorders based on the hard disk unit are emerging (refer to, “Devices Capable of Recording TV Programs on Hard Disk Drives Are Coming to Market One After Another” and NIKKEI Electronics, No. 727, pp. 41-46, 1998, “HDD-Based Digital Recording Technologies Are At Last Coming into Home” for example).\nIn the United States for example, in order to eliminate audience inconveniences such as that audiences cannot view desired broadcast programs because of unexpected telephone calls or visitors, commercially available are devices which store received broadcast content on a mass storage device such as a hard disk drive to allow audiences to view the continuation of the content later.\nHowever, unconditional storage of broadcast content by terminal audiences makes it difficult to control unauthorized content duplication activities, resulting in the violation of content copyright holders such as broadcast companies for example. Especially, the duplication and alteration of digital data and content are extremely easy, thereby exposing their copyrights to the risk. Consequently, it is desirable to limit in some way of other the use of broadcast content such as recording and reproduction on the recording/reproducing devices owned by audiences.\nFor example, limiting the recording retention period or reproducible period of the content stored by audiences to some length of time can prevent them from reusing the stored content without permission of the copyright holder; for example, the recorded content can be viewed within one hour after its broadcast time. The content recording/reproducing device may be adapted not to hold the received content which has passed a predetermined period after its broadcast or reception time (more specifically, the content is deleted from the hard disk drive of this device) or not to reproduce the expired content (more specifically, the content is retrieved from the hard disk drive).\nThus, setting a recording retention period or a reproducible period to content can technologically control the unauthorized use of content. However, if audiences do not understand the mechanism of the limited use of content or do not know the valid remaining retention or reproducible period of content, it may invite misunderstanding of audiences that the recording/reproducing device has failed or make them uncomfortable because the recording/reproducing device rejects the resumption of the reproduction of the content disrupted by an unexpected telephone call or visitor for example."}
{"text": "Graphics processing typically involves coordination of two processors, a central processing unit (CPU) and a graphics processing unit (GPU). The GPU is a specialized electronic circuit designed to accelerate the creation of images in a frame buffer intended for output to a display. GPUs are used in embedded systems, mobile phones, personal computers, tablet computers, portable game devices, workstations, and game consoles. A GPU is typically designed to be efficient at manipulating computer graphics. GPUs often have a highly parallel processing architecture that makes the GPU more effective than a general-purpose CPU for algorithms where processing of large blocks of data is done in parallel.\nThe CPU may send the GPU instructions, commonly referred to as draw commands, that instruct the GPU to implement a particular graphics processing task, e.g. render a particular texture that has changed with respect to a previous frame in an image. These draw commands may be coordinated by the CPU with a graphics application programming interface (API) in order to issue graphics rendering commands that correspond to the state of the particular application's virtual environment.\nIn order to render textures for a particular program, a GPU may perform a series of processing tasks in a “graphics pipeline” to translate the visuals in the virtual environment into images that can be rendered onto a display. A typical graphics pipeline may include performing certain rendering or shading operations on virtual objects in the virtual space, transformation and rasterization of the virtual objects in the scene to produce pixel data suitable for output display, and additional rendering tasks on the pixels (or fragments) before outputting the rendered image on a display.\nVirtual objects of an image are often described in virtual space in terms of shapes known as primitives, which together make the shapes of the objects in the virtual scene. For example, objects in a three-dimensional virtual world to be rendered may be reduced to a series of distinct triangle primitives having vertices defined in terms of their coordinates in three-dimensional space, whereby these polygons make up the surfaces of the objects. Each polygon may have an associated index that can be used by the graphics processing system to distinguish a given polygon from other polygons. Likewise, each vertex may have an associated index that can be used to distinguish a given vertex from other vertices. A graphics pipeline may perform certain operations on these primitives to produce visuals for the virtual scene and transform this data into a two-dimensional format suitable for reproduction by the pixels of the display. The term graphics primitive information (or simply “primitive information”), as used herein, is used to refer to data representative of a graphics primitive. Such data includes, but is not limited to, vertex information (e.g., data representing vertex positions or vertex indices) and polygon information, e.g., polygon indices and information that associates particular vertices with particular polygons.\nAs part of the graphics pipeline, the GPU may perform rendering tasks by implementing programs commonly known as shaders. A typical graphics pipeline may include vertex shaders, which may manipulate certain properties of the primitives on a per-vertex basis, as well as pixel shaders (also known as “fragment shaders”), which operate downstream from the vertex shaders in the graphics pipeline and may manipulate certain values on a per-pixel basis before transmitting the pixel data to a display. The fragment shaders may manipulate values relevant to applying textures to primitives. The pipeline may also include other shaders at various stages in the pipeline, such as geometry shaders that use the output of the vertex shaders to generate a new set of primitives, as well as compute shaders (CS) which may be implemented by a GPU to perform certain other general computational tasks.\nGraphical display devices having a wide field of view (FOV) have been developed. Such devices include head mounted display (HMD) devices. In an HMD device, a small display device is worn on a user's head. The display device has a display optic in front of one eye (monocular HMD) or each eye (binocular HMD). An HMD device typically includes sensors that can sense the orientation of the device and change the scene shown by the display optics as the user's head moves. Conventionally, most stages of rendering scenes for wide FOV displays are performed by planar rendering where all parts of the screen have the same number of pixels per unit area.\nTo provide a realistic experience it is desirable for the graphics presented by a wide FOV display device to be of high quality and efficiently rendered.\nIt is within this context that the present disclosure arises."}
{"text": "Diagnostic chips, and biochemical reaction chips using solid phase carriers such as protein microarrays are very useful in the field of research and development and the clinical field. For example, in terms of a protein microarray or the like, a protein solution is brought into contact with a glass substrate to immobilize the protein on the substrate, and using the physiological activity of the immobilized protein, diagnosis or analysis is carried out. In this case, it is required that the protein can be immobilized at a high density and that the substance to be reacted with the immobilized protein not adsorb nonspecifically to the carrier.\nTo eliminate such non-specific adsorption, for example, in Proteomics, (2003) Vol. 3, pp. 254-264, a method wherein a protein is directly adsorbed and immobilized via a polylysine coating layer or an amino group layer and a method wherein a protein is immobilized by covalent bonding via an active ester layer formed by a multistep reaction are described.\nFurther, in Japanese Patent Application Laid-Open (JP-A) No. 2005-201901, a method is described wherein a substrate is directly carboxylated using a particular silane coupling agent and thereby activated to allow a protein to be covalently bound to the substrate. In cases where a silane coupling agent such as the one disclosed in JP-A 2005-201901 is used, in general, after coating a substrate with the silane coupling agent, the coating is solidified by a baking treatment at a temperature of not less than 100° C. to form a three dimensional network structure.\nIn such a method, although it is possible to immobilize a protein stably at a relatively high density, the density at which the protein is immobilized is not sufficient, and non-specific adsorption of the substrate to be reacted with the immobilized protein increases.\nFurther, in order to increase the amount of immobilization of a physically active ligand such as a protein, a solid phase carrier having a porous structure such as a non-woven fabric structure of fibers or a particle-packed structure is used.\nFor example, in JP 6-148190 A, a biochemical assay is carried out by allowing an acid anhydride to be covalently bound to a porous membrane. Further, in JP 8-506902 A, a biochemical assay is carried out by immobilizing an acid anhydride-containing polymer and a protein on a porous solid phase including glass-fiber paper. Further, in WO 2006/058237, an acid anhydride-containing polymer and a protein are immobilized, and when a biochemical assay is carried out, a blocking agent is used, thereby suppressing non-specific adsorption.    [Patent Document 1] JP-A 2005-201901    [Patent Document 2] JP-A 6-148190    [Patent Document 3] JP-A 8-506902    [Patent Document 4] WO 2006/058237    [Non-patent Document 1] Proteomics, (2003) Vol. 3, pp. 254-264"}
{"text": "A display for displaying images is one of light emitting devices indispensable in the modern life and takes various forms such as a so-called TV monitor, a liquid crystal display developed rapidly in recent years, an organic EL display expected to be developed, and the like depending on the usage. Among the above, the organic EL display is best noted as a next generation flat panel display device.\nA light emission mechanism of an electroluminescent device constituting the organic EL display is such that a luminescent layer made from a luminescent composition is provided between electrodes so that electrons injected from a cathode are recombined with holes injected from an anode at the recombination center of the luminescent layer to form molecular excitons when a current is supplied and photons discharged when the molecular excitons return to the ground state are used for the light emission. Accordingly, one of preconditions for manufacturing a light emitting device of good efficiency is to inject the holes and the electrons efficiently into the luminescent layer made of an organic thin film or the like.\nUnder typical electroluminescent device operation conditions, a current of about 100 mA/cm2 is injected into the organic thin film inherently having a high electrical resistance. In order to realize such high density current injection, it is necessary to keep the sizes of a barrier against the holes injected from the anode and a barrier against the electrons injected from the cathode as small as possible. That is to say, it is necessary to use a metal having a small work function for the cathode and to select an electrode having a large work function for the anode. By selecting various metals and alloys for the cathode, it is practically possible to control the work function at will. In contrast, since transparency is required of the anode in general electroluminescent devices, the material to be used for the anode is limited to transparent conductive oxides under the current situation, and there is no alternative but to select some oxide conductive films such as an indium-tin oxide (hereinafter abbreviated to ITO) film in view of stability, transparency, resistivity, and the like at present. A work function of the ITO film can be changed to a certain degree by a history of film formation and a surface treatment, but such process has its limit. This inhibits the reduction in the hole injection barrier.\nAs one of methods to reduce the barrier against hole injection from the ITO cathode, an insertion of a buffer layer on the ITO film is known. By optimizing an ionization potential of the buffer layer, it is possible to reduce the hole injection barrier. The above-described buffer layer is called a hole injection layer. Materials which can function as the hole injection layer are generally classified into metal oxides, low molecular organic compounds, and high molecular compounds. Examples of the metal oxides are vanadium oxide, molybdenum oxide, ruthenium oxide, aluminum oxide, and the like (Non-Patent Documents 1 and 2). Examples of the low molecular organic compounds are starburst amines (Non-Patent Document 3) such as m-MTDATA, metal phthalocyanine (Patent Document 1, Non-Patent Document 4), and the like. As the high molecular compounds materials, conjugated polymers such as polyaniline (Non-Patent Document 5) and a polythiophene derivative (Non-Patent Document 6) are known. By using the above-described materials for the hole injection layer, the hole injection barrier is reduced and the holes are efficiently injected, thereby improving the efficiency and life of the electroluminescent device and reducing a driving voltage.\n[Non-Patent Document 1]\n    S. Tokito, Et al., J. Phys. D 1996, 29, 2750–2753[Non-Patent Document 2]    Tsunataka Kurosaka, et al., Shingakugiho, 1998, 98, 63–68[Non-Patent Document 3]    Y Shirota, et al., Appl. Phys. Lett. 1994, 65, 807–809[Patent Document 1]    S. A. Vanslyke, et al., U.S. Pat. No. 4,720,432[Non-Patent Document 4]    S. A. Vanslyke, et al., Appl. Phys. Lett. 1996, 69, 2160–2162[Non-Patent Document 5]    Y. Yang, et al., Appl. Phys. Lett. 1994, 64, 1245–1247[Non-Patent Document 6]    S. A. Carter, et al., Appl. Phys. Lett. 1997, 70, 2067–2069\nAmong the above-described materials, the low molecular compounds such as the starburst amines and the metal phthalocyanine, and the high molecular materials such as polyaniline are frequently used (see, for example, Patent Documents 5 and 6). Particularly, copper phthalocyanine is one of the most-frequently-used hole injection materials. This material can be obtained easily and is chemically and thermally stable. However, metal phthalocyanines have a remarkably low solubility and it is difficult to chemically modify the compounds. Further, though the materials have excellent characteristics as the hole injection material, they are colored in many cases to disadvantageously color a light emission surface itself of the device. Therefore, there is a demand for an electroluminescent device using a hole injection material having no absorption or small absorption intensity in the visible region."}
{"text": "PTEN (phosphatase and tensin homolog deleted on chromosome 10) is the protein which is encoded by the PTEN gene. The PTEN protein is found in almost all tissues in the body. PTEN is involved in the regulation of the cell cycle by preventing cells from growing and dividing too rapidly. The PTEN gene was identified as a tumor suppressor gene that is mutated in a large number of cancers at high frequency. These include cancers of the breast, prostate, endometrium, ovary, brain, skin, thyroid, lung, bladder and colon, as well as melanoma, glioblastoma. and lymphoma (Teng et al., Cancer Res. 57: 5221-5225, 1997) During tumor development, mutations and deletions of PTEN occur that inactivate its enzymatic activity leading to increased cell proliferation and reduced cell death. Genetic inactivation of PTEN occurs in glioblastoma, endometrial cancer, prostate cancer, and reduced expression is found in many other tumor types, such as lung and breast cancer.\nPTEN mutations also cause a variety of inherited predispositions to cancer. For example, Cowden syndrome, Bannayan-Riley-Ruvalcaba syndrome, Proteus syndrome and Proteus-like syndrome. Together, the disorders caused by PTEN mutations are called PTEN hamartoma tumor syndromes, or PHTS. Mutations responsible for these syndromes cause the resulting protein to be nonfunctional or absent. The defective protein allows the cell to divide in an uncontrolled way and prevents damaged cells from dying, which can lead to the growth of tumors.\nThe prevalence of PTEN mutations in cancer underscores the significance of finding methods and compositions to treat, diagnose, and prevent such diseases. However, perhaps because of the many cellular biochemical reactions in which PTEN phosphatase activity is implicated in, or because of the nature of tumor suppressor genes, treatments that directly target PTEN activity or expression have had limited success. Thus, a need exists to develop new methods and compositions to identify patients suitable for treatment, treatments, diagnosis and prognosis of cancer and related diseases in which PTEN is mutated or lost."}
{"text": "Conventionally, for example, in automobile electric components such as automobile lamps and ECUs (Electrical Control Units), OA (office automation) equipment, household electric appliances, and medical devices, a housing containing an electric component or a control board is provided with an opening for the purpose of reducing pressure variation in the housing caused by temperature change or allowing replacement of air in the housing, and a ventilation member is attached to the opening. The ventilation member ensures ventilation between the inside and outside of the housing and also prevents foreign matters such as dust and water from entering the housing.\nFor example, Patent Literature 1 discloses a ventilation member 100 as shown in FIGS. 9A and 9B. The ventilation member 100 includes: a water-proof gas-permeable membrane 120 for covering an opening 151 of a housing 150; a support 110 supporting the water-proof gas-permeable membrane 120 and configured to be fixed to the housing 150; and a cover component 130 covering the water-proof gas-permeable membrane 120.\nThe support 110 has a tubular shape as a whole, and includes: a flange portion 112 to which the water-proof gas-permeable membrane 120 is adhered; and an insertion portion 111 projecting from the flange portion 112 and configured to be inserted into the opening 151 of the housing 150. A sealing member 140 is mounted around the root portion of the insertion portion 111. The sealing member 140 is pressed against the housing 150 by the flange portion 112, so that a gap between the ventilation member 100 and the housing 150 is sealed.\nThe cover component 130 includes: a principal wall 131 facing the water-proof gas-permeable membrane 120; and a peripheral wall 132 extending from the peripheral portion of the principal wall 131 to the vicinity of a surface of the housing 150. In addition, the principal wall 131 is provided with through holes 133 for allowing a space facing the water-proof gas-permeable membrane 120 to communicate with the outside.\nIn addition, Patent Literature 2 discloses a ventilation structure 300 in which partition walls 152 are provided on the housing 150 so as to surround a ventilation member 200 as shown in FIGS. 10A and 10B. The partition walls 152 are arc-shaped walls spaced from each other, and have a height greater than that of the ventilation member 200."}
{"text": "Field\nEmbodiments described herein generally relate to methods and devices for self-calibration of an actuator. More specifically, embodiments described herein relate to methods and devices for self-calibration of a linear actuator configured to control a spring return valve with variable stroke.\nBackground\nCalibration of an actuator ensures proper operation and optimal performance. One approach is to use mechanical switches to detect stroke limits. Due to the mechanical contact, the accuracy, stability, and reliability of such systems may be significantly compromised. Furthermore, the mechanical switches take up additional space.\nAnother approach is to use resistors as current sensors for signal detection of motor power lines. However, due to variation in current flow, the accuracy and reliability of such sensor systems may be adversely affected. Furthermore, the sensing resistor results in poor system efficiency, resulting in using a larger actuator for the same application because the resistance must be large enough for sufficient signal sensing."}
{"text": "This invention relates to new chemical compositions, and more particularly to double chromates having the formula MM'(CrO.sub.4).sub.2 or MM'(Cr.sub.2 O.sub.7).sub.2 or MM' (CrO.sub.x).sub.2 where M and M' represent different metals and where x is less than 4 and greater than or equal to 3.7, to methods of making them, and to their use as positive active materials in batteries.\nFrom the disclosure of French Pat. No. 2 110 765, a lithium battery is known having a nonaqueous electrolyte and positive active material constituting silver chromate having the formula Ag.sub.2 CrO.sub.4. One of the great interests in this type of battery is that the reduction in metallic silver occurs, with sufficiently small currents, at a level of about 3 volts, with good yield and stability, so that such a battery can be used in applications such as clocks or cardio-stimulators. This active material has given fully satisfactory results, but it is relatively expensive.\nOther chromates have been proposed as the active positive material for batteries, for example, chromates of copper, of iron, of cobalt, of nickel, of mercury, of thallium, and of bismuth, but aside from the fact that the existence of certain of these chromates is doubtful, their stability and discharge voltage would be inferior to the above-mentioned silver chromate."}
{"text": "As an optotype presenting apparatus of a vision testing system, an apparatus arranged to present an optotype or chart on a color liquid crystal display (LCD) has recently been proposed as well as a conventional apparatus arranged to present an optotype on a screen by projection. Furthermore, an apparatus with the LCD for a binocular balance test has been proposed in which two LCDs placed one before the other and two polarization filters having polarization axes perpendicular to each other. In this apparatus, as with the conventional screen projection type apparatus, two polarization filters whose polarization axes are perpendicular to each other are disposed in front of both eyes of an examinee to present an optotype for right eye and an optotype for left eye to both eyes of the examinee (see U.S. Pat. No. 5,331,358 (JP5-130975 (1993)).\nHowever, such configuration having two LCDs and two polarization filters is complicated and expensive and leads to an increase in apparatus size."}
{"text": "A borehole is typically drilled using a drill bit which is attached to an end of a drill string. Rotary drilling is performed by rotating the drill bit. The drill bit may be rotated by rotating the drill string, by rotating the drill bit with a downhole drilling motor, or in some other manner.\nA roller cone drill bit includes cones which rotate as the drill bit is rotated. Teeth which are positioned on the cones roll along the bottom of the borehole as the cones rotate. The teeth impact the bottom of the borehole as they roll and thereby crush and disintegrate rock in order to advance the borehole.\nA fixed cutter drill bit typically includes no moving parts, but includes cutters which are attached to the body of the drill bit and which rotate with the drill bit as the drill bit is rotated. The cutters scrape the borehole as the drill bit rotates, thereby shearing rock in order to advance the borehole.\nA cutter on a fixed cutter drill bit is typically comprised of a cutter element, such as an “abrasive” or “superabrasive” cutter element, which performs the shearing action. An abrasive cutter element may be comprised of tungsten carbide, another carbide material, ceramic and/or some other material. A superabrasive cutter element may be comprised of natural diamond, a synthetic diamond material such as polycrystalline diamond compact (PDC) or thermally stable diamond (TSP), or may be comprised of some other material such as cubic boron compact or diamond grit impregnated substances.\nA cutter on a fixed cutter drill bit may be further comprised of a substrate to which the cutter element may be affixed. For example, a PDC or TSP cutter element may be comprised of a disc or cylinder shaped “diamond table” which may be affixed to a substrate such as tungsten carbide in order to provide the complete cutter. The diamond table typically comprises a substantially flat and circular cutting face which contacts the borehole in order to perform the shearing action.\nA PDC or TSP cutter element may typically be affixed to a substrate by applying high temperature and high pressure to the cutter element and substrate in the presence of a catalyst so that the materials of the cutter element and the substrate bond with each other.\nFixed cutter drill bits are therefore typically comprised of a bit body and a plurality of cutters which are attached to the bit body. The bit body is typically constructed of steel or of a matrix containing an erosion resistant material such as tungsten carbide. The cutters are typically attached to the bit body by an adhesive or by brazing. The cutters may be received in cutter pockets in the bit body in order to facilitate the attachment of the cutters to the bit body.\nThe bit body and the cutters are configured to provide an overall design for the drill bit, having regard to considerations such as rate of penetration of the drill bit, drill bit stability, drill bit steerability, drill bit durability and hydraulic performance of the drill bit.\nFor example, the bit body typically includes a plurality of blades to which the cutters are attached and between which fluids and cuttings may pass. Because the cutters are typically attached to the blades of the drill bit, increasing the number of blades on a fixed cutter drill bit will generally increase the number of cutters which may be attached to the bit body, thereby increasing the “cutter count” and the “cutter density” on the drill bit.\nGenerally, the rate of penetration which can be achieved by a fixed cutter drill bit is inversely proportional to the number of blades and cutters which are included in the drill bit. In other words, the greater the number of blades and the greater the number of cutters, the lower the rate of penetration which may be expected from the drill bit.\nGenerally, the durability of the drill bit is proportional to the number of blades and cutters which are included in the drill bit. In other words, the greater the number of blades and the greater the number of cutters, the longer the drill bit may be expected to function without experiencing excessive wear.\nGenerally, the hydraulic performance of the drill bit is inversely proportional to the number of blades which are included in the drill bit. In other words, the greater the number of blades, the less area which is available between the blades for the passage of fluids and cuttings, and the more resistance which is provided to the passage of fluids and cuttings past the drill bit.\nAs a result, the design of a fixed cutter drill bit typically represents a compromise amongst the rate of penetration, stability, steerability, durability, and hydraulic performance which can be achieved with the drill bit. Various design strategies have been proposed for achieving an appropriate balance of these considerations.\nU.S. Pat. No. 6,283,233 (Lamine et al) describes a drilling and/or coring tool which includes PDC cutting elements and/or secondary cutting elements and at least one associated cutting element which is situated behind at least one of the PDC or secondary cutting elements, wherein the associated cutting element is “hidden” behind the PDC or secondary cutting element and is unused unless or until the PDC or secondary cutting element with which it is associated wears down, is torn away, or is broken.\nU.S. Pat. App. Pub. No. US 2006/0070771 A1 (McClain et al), U.S. Pat. App. Pub. No. US 2007/0079995 A1 (McClain et al), and U.S. Pat. App. Pub. No. US 2008/0149393 A1 (McClain et al) all describe a drill bit for drilling through a casing bit which is disposed at the end of a casing. The drill bit includes a first type of cutting element and a second type of cutting element. The first type of cutting element is comprised of a superabrasive material and the second type of cutting element may be comprised of either a superabrasive material or an abrasive material. The second type of cutting element is positioned behind the first type of cutting element but exhibits a “relatively greater exposure” than the first type of cutting element so as to engage the interior of the casing bit and drill through the casing bit. The second type of cutting element then wears quickly upon engagement with the subterranean formation, after which the first type of cutting element continues to drill through the subterranean formation.\nU.S. Pat. App. Pub. No. US 2007/0199739 A1 (Schwefe et al) describes a cutter insert for a fixed cutter drill bit which may be used to secure a backup cutter in a recess behind a primary cutter on the drill bit. The backup cutter may be configured to be underexposed, overexposed or to have a substantially equal exposure relative to the primary cutter.\nU.S. Pat. App. Pub. No. US 2008/0179106 A1 (Gavia et al) and U.S. Pat. App. Pub. No. US 2008/0179108 A1 (McClain et al) both describe a rotary drag bit which includes a primary cutter and at least two additional cutters which are positioned on a single blade of the drill bit and which are configured relative to each other. In particular, the additional cutters are configured to follow the primary cutter.\nU.S. Pat. App. Pub. No. US 2008/0179107 A1 (Doster) describes a rotary drag bit which includes a plurality of blades and at least one split cutter set. The split cutter set includes a plurality of cutters, where at least two of the cutters are primary and/or kerfing cutters located on different blades of the bit, and where at least one of the cutters is a backup cutter. The cutters in the split cutter set all follow substantially a common cutting path upon rotation of the bit body about its central axis.\nPCT International Publication No. WO 2008/095005 A1 (Chen et al) describes a rotary drill bit with cutting elements which are operable to control the depth of cut and rate of penetration during drilling of a wellbore. The cutting elements may be arranged in sets of a primary cutting element and an associated secondary cutting element, wherein the secondary cutting element is disposed in a leading position relative to the primary cutting element, and wherein the cutting face of the primary cutting element is exposed a greater distance from the bit face profile than the cutting face of the secondary cutting element. The sets of cutting elements may also be comprised of a “protector” which is operable to control the depth of the cut of the cutting elements.\nThere remains a need for fixed cutter drill bits which facilitate reasonable compromises with respect to the rate of penetration, stability, steerability, durability, and hydraulic performance which can be achieved with the drill bit."}
{"text": "Integrated circuits (ICs), such as, ultra-large scale integrated (ULSI) circuits, can include as many as one million transistors or more. The ULSI circuit can include complementary metal oxide semiconductor (CMOS) field effect transistors (FETS). The transistors can include semiconductor gates disposed between drain and source regions. The drain and source regions are typically heavily doped with a P-type dopant (boron) or an N-type dopant (phosphorous).\nThe drain and source regions generally include a thin extension that is disposed partially underneath the gate to enhance the transistor performance. Shallow source and drain extensions help to achieve immunity to short-channel effects which degrade transistor performance for both N-channel and P-channel transistors. Short-channel effects can cause threshold voltage roll-off and drain-inducted barrier-lowering. Shallow source and drain extensions and, hence, controlling short-channel effects, are particularly important as transistors become smaller.\nConventional techniques utilize a double implant process to form shallow source and drain extensions. According to the conventional process, the source and drain extensions are formed by providing a transistor gate structure without sidewall spacers on a top surface of a silicon substrate. The silicon substrate is doped on both sides of the gate structure via a conventional doping process, such as, a diffusion process or an ion implantation process. Without the sidewall spacers, the doping process introduces dopants into a thin region (i.e., just below the top surface of the substrate) to form the drain and source extensions as well as to partially form the drain and source regions.\nAfter the drain and source extensions are formed, silicon dioxide spacers, which abut lateral sides of the gate structure, are provided over the source and drain extensions. The substrate is doped a second time to form the deeper source and drain regions. The source and drain extensions are not further doped due to the blocking capability of the silicon dioxide spacer.\nAs transistors disposed on integrated circuits (ICs) become smaller, transistors with shallow and ultra-shallow source/drain extensions have become more difficult to manufacture. For example, smaller transistors should have ultra-shallow source and drain extensions with less than 30 nanometer (nm) junction depth. Forming source and drain extensions with junction depths of less than 30 nm is very difficult using conventional fabrication techniques. Conventional ion implantation and diffusion-doping techniques make transistors on the IC susceptible to short-channeling effects, which result in a dopant profile tail distribution that extends deep into the substrate. Also, conventional ion implantation techniques have difficulty maintaining shallow source and drain extensions because point defects generated in the bulk semiconductor substrate during ion implantation can cause the dopant to more easily diffuse (transient enhanced diffusion, TED). The diffusion often extends the source and drain extension vertically into the bulk semiconductor substrate.\nThus, there is a need for a method of manufacturing ultra-shallow source and drain extensions that is not susceptible to transient enhanced diffusion. Further still, there is a need for transistors that have ultra-shallow junction source and drain extensions. Even further still, there is a need for an efficient method of manufacturing source and drain extensions."}
{"text": "The present invention generally relates to manufacturing, remanufacturing or repairing replaceable imaging components, and more particularly to apparatus and techniques for providing a drive gear or a non-drive gear for a drum or roller, such as an organic photo conductor (OPC) drum, for example, of a replaceable imaging cartridge adapted for holding marking material, such as toner.\nIn the imaging industry, there is a growing market for the remanufacture and refurbishing of various types of replaceable imaging cartridges such as toner cartridges, drum cartridges, inkjet cartridges, and the like. These imaging cartridges are used in imaging devices such as laser printers, xerographic copiers, inkjet printers, facsimile machines and the like, for example. Imaging cartridges, once spent, are unusable for their originally intended purpose. Without a refurbishing process these cartridges would simply be discarded, even though the cartridge itself may still have potential life. As a result, techniques have been developed specifically to address this issue. These processes may entail, for example, the disassembly of the various structures of the cartridge, replacing toner or ink, cleaning, adjusting or replacing any worn components and reassembling the imaging cartridge.\nLaser printer toner cartridges are typically composed of two portions. One of these sections is the waste bin assembly which houses the OPC drum. The OPC may include a drive gear which engages with a printer drive member. During the remanufacturing of a laser printer toner cartridge, the OPC drum may need to be replaced due to the wear or damage of the OPC drum. The replacement OPC drum may include a replacement drive gear attached to one end of the replacement OPC drum and a replacement non-drive hub attached to the opposing end of the replacement OPC drum. The present invention provides for an improved replacement drive gear and a replacement non-drive hub."}
{"text": "One of the most CPU intensive activities associated with performing network protocol processing is the need to copy incoming network data from an initial landing point in system memory to a final destination in application memory. This copying is necessary because received network data cannot generally be moved to the final destination until the associated packets are: A) analyzed to ensure that they are free of errors, B) analyzed to determine which connection they are associated with, and C) analyzed to determine where, within a stream of data, they belong. Until recently, these steps had to be performed by the host protocol stack. With the introduction of the intelligent network interface device (as disclosed in U.S. patent application Ser. Nos. 09/464,283, 09/439,603, 09/067,544, and U.S. Provisional Application Ser. No. 60/061,809), these steps may now be performed before the packets are delivered to the host protocol stack.\nEven with such steps accomplished by an intelligent network interface device, there is another problem to be addressed to reduce or eliminate data copying, and that is obtaining the address of the destination in memory and passing that address to the network interface device. Obtaining this address is often difficult because many network applications are written in such a way that they will not provide the address of the final destination until notified that data for the connection has arrived (with the use of the “select( )” routine, for example). Other attempts to obtain this address involve the modification of existing applications. One such example is the Internet Engineering Task Force (IETF) Remote DMA (RDMA) proposal, which requires that existing protocols such as NFS, CIFS, and HTTP be modified to include addressing information in the protocol headers. A solution is desired that does not require the modification of existing applications or protocols."}
{"text": "1. Field\nA spike fixing member and a cloth treating apparatus having the same are disclosed herein.\n2. Background\nRecently, various kinds of cloth treating apparatuses other than washing machines have been developed and used. Representatively, there is a cloth treating apparatus having various functions of drying washed clothes, removing wrinkles in clothes by applying hot air or water, removing a smell remaining on the clothes, and so on. The cloth treating apparatus may be variously referred to as a cloth managing apparatus, a refresher, or a styler.\nIn Korean Patent Publication No. 10-2015-0085728, there is disclosed a cloth treating apparatus having a structure in which dried air heated while passing through a condenser by driving of a heat pump including a compressor, the condenser, an evaporator, and an expander may be supplied to a cloth accommodation space and dries clothes, and moist air in the cloth accommodation space may be dehumidified and cooled while passing through the evaporator and then circulated while passing through the condenser. Condensate water generated while the air passes through the evaporator may be collected in a container. In Korean Patent Publication No. 10-2009-0102400, there is disclosed a cloth treating apparatus in which a steam generating unit and a steam spraying unit for spraying steam into a cloth accommodation space are provided to remove wrinkles of accommodated clothes and to enhance deodorization performance.\nThe cloth treating apparatus having such a structure may be provided at an indoor space. When the cloth treating apparatus is provided at the indoor space, a horizontal control thereof is essential for a normal operation of the cloth treating apparatus. To this end, a leg which adjusts a height may be provided at a lower surface of a cabinet of the steam generating unit.\nA representative structure of the leg is disclosed in Korean Patent Publication No. 10-2015-0042948. The structure of the leg may adjust the height and may have a vibration preventing member to reduce the vibration.\nHowever, when the cloth treating apparatus is provided an indoor space on which carpet is spread, it may be difficult to stably fix the cloth treating apparatus. If it is impossible to stably fix the cloth treating apparatus, vibration may occur when the cloth treating apparatus is operated. The above references are incorporated by reference in their entirety herein where appropriate for appropriate teachings of additional or alternative details, features and/or technical background."}
{"text": "1. Field of the Invention\nThe present invention concerns a system and process for ensuring the provenance of passengers who intend to board a vehicle at a transportation facility such as, e.g., a commercial airport.\n2. Discussion of the Known Art\nCertain public transportation systems or facilities must ensure the provenance of each passenger before allowing him or her to board a vehicle such as a commercial aircraft at the facility. As used herein, “provenance” implies (i) knowledge that a given person who has checked into a transportation facility is the same person who later appears at a departure area (gate) of the facility or who actually enters the vehicle, and (ii) knowledge of the person with respect to (a) the ability of him or her to pass successfully through immigration or other checks on arrival, and (b) a security assessment based, for example, on a background check in order to quantify a risk the passenger might pose to other individuals (passengers or otherwise) if allowed to board the vehicle. Knowledge of both the mentioned items (i) and (ii) will therefore allow a determination to be made of the passenger's overall fitness for travel through the facility.\nU.S. Pat. No. 6,119,096 (Sep. 12, 2000) relates to a system and method for automated aircraft boarding and passenger billing, wherein passenger check-in and boarding is determined by the use of an iris recognition system in conjunction with associated passenger accounts. The patent does not address privacy issues concerning the anonymity of the biometric record, or certain instances in actual deployments wherein data attempted to be acquired from the iris is insufficient alone for recognition of a particular passenger (e.g., a blind person whose eyelids are closed) without other information. In addition, the patent is not concerned with the fitness of prospective passengers for boarding and traveling on an aircraft, nor does it disclose particular billing methods that may be optimal with respect to the transportation carrier or organization."}
{"text": "This invention relates to a process for thermal upgrading of low rank coal, for instance brown coal and other carbonaceous material, for instance peat, lignite, wood or carbonaceous waste. In this specification such other carbonaceous material will be deemed to be included in the term \"coal\".\nLow rank coal, in the least favorable case can contain up to 90%w water, but more generally has a water content of 40 to 70%. The water is present in various forms, namely, chemically bound, gel, and absorbed water. This water needs to be removed from the coal not only for efficient transport, but also in order to increase its calorific value (and thus its worth) and to improve its combustion properties. This water can be removed either partially by thermal drying, in which the surface and absorbed water is evaporated, or completely by thermal upgrading. In general, water removed by thermal drying will over a period of time be reabsorbed, but water removed by thermal upgrading will not, due to chemical changes which take place in the coal.\nThermal upgrading of coal can be divided into two distinct stages, namely, dewatering, and decarboxylation (removal of oxygen-containing groups). Dewatering generally takes place at about 200.degree. C. and results in a shedding of an important part of the bound water. On separating the water from the coal, there is not a tendency for it to be reabsorbed Decarboxylation takes place at a rather higher temperature, particularly above 300.degree. C. and results in a further structural change in the coal. The coal significantly becomes hydrophobic due to the removal of the oxygen-rich polar groups which are responsible for the hydrophilic properties of the coal. For instance, confirmation of the structural change of the coal is found from petrographic analysis, from which it is seen that the vitrinitic reflectance for a typical brown coal as a result of treatment above 300.degree. C. has improved from ca. 0.35 before treatment to 0.70, which is equivalent to that of a subbituminous coal.\nWhile for some types of coal thermal drying can be used with advantage, where the coal has a high water content it is less efficient in that the latent heat of vaporization of the water has to be supplied. Furthermore, there is practically no change in structure of the coal, so that the resulting dried product remains hydrophilic, and in the absence of special precautions will thus reabsorb much water in a humid environment. And, the calorific value may not be sufficiently increased to command a price which can justify the expense of transport over any considerable distance. For example, prior processes for dewatering wet, solid carbonaceous fuel materials have been described in patents such as the following: U.K. Pat. No. 844,556 describes removing water and bitumen from oil shale or peat or the like, by heating the moist material in oil, first at a temperature at which the water evaporates, then at a temperature at which the bitumen forms and evaporates. U.S. Pat. No. 3,552,031 describes heating brown coal as received from the mines at from about 100.degree. to 300.degree. C. at a pressure which prevents water vaporization but is less than 5000 psi, mechanically separating a fines-containing aqueous liquid which is subsequently separated into waste liquid and a sludge which is recycled into the coal to be treated. U.S. Pat. No. 3,992,784 by E. Verschuur, B. P. Ter Meulen, T. V. Herwijnen and J. Boom describes heating a slurry of brown coal and water at a temperature above 150.degree. C. and a pressure which prevents the vaporization of water, mechanically separating this slurry into coal, water and wet fines and recycling at least some of the water to preheat the incoming coal."}
{"text": "1. Technical Field of the Invention\nThe present invention generally relates to integrated circuits. More particularly, and not by way of any limitation, the present invention is directed to a system and method for measuring fault coverage in an integrated circuit.\n2. Description of Related Art\nThe evaluation of the reliability and quality of a digital integrated circuit (IC) comprises several distinct testing phases including a verification phase, a parametric testing phase and a defect testing phase. The verification testing phase is the initial phase in which the first prototype chips are tested to ensure that they match their functional specification, that is, to verify the correctness of the design. The verification testing phase checks that all design rules are adhered to, from layout to electrical parameters. The parametric testing phase ensures components meet design specification for delays, voltages, power, etc.\nThe defect testing phase ensures that only defect-free production chips are packaged and shipped. The defect testing phase involves testing for a relatively large number of different physical defects that could be present. These physical defects may be classified by their logical fault effect on the functionality of the circuit. For example, permanent faults are physical defects which exist long enough to be observed. Temporary faults are physical defects which appear and disappear in short intervals of time. Delay faults are physical defects which affect the operating speed of the circuit. Sequential faults are physical defects which cause a combinational circuit to behave like a sequential circuit and occur in only certain technologies (e.g. CMOS).\nOne of the most practical methods for defect testing is to employ a fault model of the physical defects that can occur in the IC at a high level of abstraction, typically the logic level, and then develop algorithms, commonly referred to as Functional Test Vector (FTV) sets, for the modeled faults. Depending on the fault model and quality of the functional test vector set, the functional test vector set may cover a high percentage of the actual physical defects. The percent of all possible internal faults in an IC that are observable by the FTV set employed is referred to as the fault coverage. In general, the more robust and reliable fault models and associated FTV sets provide greater fault coverage.\nProcesses for automated fault testing of large scale integrated circuits commonly employ what is known as a “stuck-at” fault model to emulate permanent faults that may occur during fabrication of the integrated circuit under test. In this model, the circuit description is modified or otherwise rendered to correspond to a stuck-at fault state, i.e., a continuous logic state of 0 or 1. In a process known as fault simulation, the test vector sets developed based on the stuck-at fault model are applied to the IC and the value of the corresponding response is compared to the expected response value. After many faults are simulated and multiple test vectors run, an indication of the fault coverage of the applied test vector set is provided. If the fault coverage is unacceptably low, the test vector set may be modified to exercise portions of the circuit where undetectable faults lie.\nTypically, fault coverage is measured via software fault simulations. In software fault simulations, a model of the IC is seeded with stuck-at faults and the functional test vector sets are run. This simulation is repeated thousands of times, seeding a different fault each time, to determine fault coverage.\nIt has been found, however, that as the size and complexity of ICs continues to grow, the software fault simulations do not provide fast and effective measurements of fault coverage. Moreover, software simulations provide a theoretical, not empirical, indication of fault coverage. Therefore, a need exists for a system and method for measuring fault coverage that is robust enough for today's production cycles that include ICs of ever-growing size and complexity. Additionally, a need exists for a system and method for measuring fault coverage that provides empirical data."}
{"text": "The present invention relates to a display and image capture apparatus which enables eye contact between callers in a videophone, video teleconferencing or similar system.\nWith recent developments of telecommunication technology such as video coding techniques, a two-way visual telecommunication system which links remote places through a video and audio communication network, such as a videophone or video teleconferencing system, has rapidly come into wide use. It is now being expected that the two-way visual telecommunication will be able to offer images full of a sense of reality by the expansion of an integrated services digital network (ISDN) and broadening of the band of the transmission network therefor. FIG. 1 shows a display and image capture apparatus which is now employed in the two-way visual telecommunication. A video camera 3 is mounted on the top or one side of a CRT display or like display 1, the image of a user M captured by the video camera 1 is transmitted to a similar display and image capture apparatus (not shown) at a remote place and displayed on its display. In this instance, however, the user M of such an apparatus is usually looking at the display 1 on which the person he is talking to is being displayed, and hence he will not turn his gaze on the video camera 3. This results in a boring conversation between callers who do not look each other in the eyes, although their images can be transmitted with a sense of reality.\nMethods which employ a half-transparent mirror (which may be referred to simply as HM) for enabling eye contact between callers are proposed in Japanese Patent Application Laid Open Nos. 269128/87 and 11082/90, for example. FIG. 2 shows the basic configuration for such methods. A half-transparent mirror 2 is disposed in front of the display 1. An image from a subject, i.e. the image of the user M is reflected by the half-transparent mirror 2 and is captured by the video camera 3 disposed above the half-transparent mirror 2. Consequently, the image-thus captured is equivalent to an image which is captured from the position of the display 1, and since the user's (i.e. talker's) eye is being directed toward the display 1, the eyes of the talker M and the listener being displayed on the display screen will meet. With such a configuration as depicted in FIG. 2, however, the size of the half-transparent mirror 2 increases with an increase in the size of the display screen of the display 1. This will also increase the depth D of the display and image capture apparatus, resulting in the apparatus inevitably becoming bulky.\nIn Japanese Patent Public Application Laid open No. 37890/90 there is proposed, as another method for enabling eye contact between callers, a configuration in which a mirror or half-transparent mirror having a number of parallel slits cut therein is disposed aslant in the same fashion as in the case of FIG. 2, in place of using the continuous half-transparent mirror. This is also defective in that the overall depth of a combined structure of the mirror and the display is large as is the case with the prior art example of FIG. 2. In this laid open application it is further stated that the depth is reduced by dividing the slit mirror into two and disposing them vertically, but in this case, since light reflected by the lower slit mirror enters the video camera through the upper slit mirror, a double image is captured; hence, this structure is impractical. In Japanese Patent Application Laid Open No. 209403/91 there is disclosed, with a view to reducing the depth of a projection TV set, a configuration wherein light from a projector provided at the lower part in the set is projected to a Fresnel mirror mounted on the ceiling of the set and the reflected light from the Fresnel mirror is projected onto the back of a screen mounted on the front of the set. Yet, this laid open patent application makes no mention of capturing the image of a person looking at the image on the TV display from the direction to meet his eye, and even if the Fresnel mirror replaces the aforementioned half-transparent mirror 2 in FIG. 1, the user M cannot see his image displayed on the display, because the Fresnel mirror is not half-transparent.\nThus, the conventional display and image capture apparatuses possess the shortcomings that no eye contact is enabled between callers or the apparatuses which enable eye contact are bulky."}
{"text": "FIG. 1 is a block diagram of a prior art system 100, according to an implementation. The system 100 includes a processor 104. A first external device 106 is physically located in the facility 102. An example of the first external device 106 is a physical security monitoring device, such as cameras, fiber optic cable, security network switches, uninterruptable power supplies and security server hard drives.\nData from the first external device 106 is transmitted to a second external device 108 over a virtual private network 110. Data from the first external device 106 is transmitted outside the facility 102. The first external device 106 is connected to a computer firewall 112 in the facility 102.\nA second external device 108 is external to the facility 102. The second external device 108 is physically located outside the facility 102. Data from the first external device 106 is transmitted to the computer firewall 112.\nThe second external device 108 includes a sentry component 114 that is operably coupled to the processor 104. The sentry component 114 performs security monitoring functions of the first external device 106, which yield security monitoring second external device 108 also includes a remediation component 116. The remediation component 116 performs remediation functions on the first external device 106 in reference and in response to the data from the first external device 106, which yields remediation results.\nBecause data from the first external device 106 is transmitted outside the facility 102, the security monitoring functions of the first external device 106 are performed in the second external device 108.\nData from the first external device 106 must be transmitted through the VPN 110 or stored outside of the facility 102 The VPN 110 presents a serious risk to the facility 102."}
{"text": "The present invention relates to a system and method for the treatment of disorders of the vasculature. More specifically, the present invention relates to a system and method for treatment of thoracic or abdominal aortic aneurysm and the like, which is a condition manifested by expansion and weakening of the aorta. Such conditions require intervention due to the severity of the sequelae, which frequently is death. Prior methods of treating aneurysms have consisted of invasive surgical methods with graft replacement within the affected vessel as a reinforcing member of the artery. However, such a procedure requires a surgical cut down to access the vessel, which in turn can result in a catastrophic rupture of the aneurysm due to the decreased external pressure from the surrounding organs and tissues, which are moved during the procedure to gain access to the vessel. Accordingly, surgical procedures have a high mortality rate due to the possibility of the rupture discussed above in addition to other factors. Other factors can include poor physical condition of the patient due to blood loss, anuria, and low blood pressure associated with the aortic abdominal aneurysm. An example of a surgical procedure is described in a book entitled Surgical Treatment of Aortic Aneurysms by Denton A. Cooley, M.D., published in 1986 by W.B. Saunders Company.\nDue to the inherent risks and complexities of surgical procedures, various attempts have been made in the development of alternative methods for deployment of grafts within aortic aneurysms. One such method is the non-invasive technique of percutaneous delivery by a catheter-based system. Such a method is described in Lawrence, Jr. et al. in “Percutaneous endovascular graft: experimental evaluation”, Radiology (May 1987). Lawrence describes therein the use of a Gianturco stent as disclosed in U.S. Pat. No. 4,580,568. The stent is used to position a Dacron® fabric graft within the vessel. The Dacron graft is compressed within the catheter and then deployed within the vessel to be treated. A similar procedure has also been described by Mirich et al. in “Percutaneously placed endovascular grafts for aortic aneurysms: feasibility study,” Radiology (March 1989). Mirich describes therein a self-expanding metallic structure covered by a nylon fabric, with said structure being anchored by barbs at the proximal and distal ends.\nOne of the primary deficiencies of the existing percutaneous devices and methods has been that the grafts and the delivery catheters used to deliver the grafts are relatively large in profile, often up to 24 French and greater, and stiff in bending. The large profile and bending stiffness makes delivery through the irregular and tortuous arteries of diseased vessels difficult and risky. In particular, the iliac arteries are often too narrow or irregular for the passage of a percutaneous device. Because of this, noninvasive percutaneous graft delivery for treatment of aortic aneurysm is not available to many patients who would benefit from it.\nWhat has been needed is an endovascular graft and delivery system for the graft which has a small outer diameter and high flexibility to facilitate percutaneous delivery in patients who require such treatment. What has also been needed is a delivery system for an endovascular graft which is simple, reliable and can accurately deploy an endovascular graft within a patient's body."}
{"text": "This invention is related to methods for forming superconductive joints; more particularly, the present invention is directed to methods for forming superconductive joints between superconductors having as many as thousands of wire filaments.\nIn whole body, nuclear magnetic resonance (NMR) imaging systems, it is necessary to provide a high strength, uniform magnetic field, typically between about 0.04 and 1.5 Tesla, or more. A desirable means for providing this field employs superconductive coil circuits operating in the persistent mode. This method for producing such magnetic fields is desirable in that continual supply of electrical energy to the coils is not required. However, such circuits must typically carry current at a high level of current density, typically 10,000 amperes/cm.sup.2, or more. To insure that such circuits do not spontaneously exhibit localized resistive (or ohmic) regions which leads to quenching of the current, it is necessary to provide sound superconductive joints on various portions of the current loop. Furthermore, because of the high levels of current employed, it is highly desirable to employ superconductive conductors having a composite structure. The composite structure typically takes a form in which a large plurality of superconductive strands or filaments are disposed within an overlaying or surrounding matrix of normal, resistive material such as copper or aluminum. Typically the copper or aluminum matrix provides an alternate, localized current path in the event that a small portion of the superconductive wire exhibits a transition to the resistive (ohmic) state. Such a composite structure is desired because of the relatively large current density desired in NMR imaging. It should also be understood that the superconductors and composite superconductors described herein are typically disposed within a coolant bath, such as liquid helium, so as to keep their temperature in the vicinity of 4.2.degree. K. to maintain the superconductive material in the zero resistance state.\nTwo methods have heretofore been used to produce superconductive joints. In a first of these methods, individual superconductive filaments, after stripping of the surrounding matrix, are individually spot welded to an intermediate niobium foil. However, this method is impractical when hundreds or thousands of filaments are to be joined. Unfortunately, modern composite superconductive conductors often employ thousands of filaments in their construction. Since such conductors are desirable in the coils for NMR imaging systems, this method of superconductive joint formation is therefore highly undesirable from a manufacturing view point. In a second method, bundles of superconductive filaments from which the surrounding matrix has been stripped, are crimped together in a copper sleeve. However, this method of joint formation has not been demonstrated to work with a strand of more than one or two hundred filaments. Thus, this method is also seen to be undesirable for the construction of the main magnetic field coils for NMR imaging systems.\nAccordingly, it is seen that it is highly desirable to be able to join together superconducting conductors having between 200 and 2,000 filaments (or more) together in a way in which a sound truly superconductive joint is formed. It is further seen that this joint should be capable of long term operation in a low temperture coolant bath and be constructed so that it does not exhibit a tendency for transition to the normal, resistive state particularly as a result of microscopic filament motion."}
{"text": "Lighting systems with a plurality of lighting units are being used today for various applications, for example for room lighting applications to create defined lighting scenes. US 2007/0258523 A1 discloses a lighting system with controllable lighting units. A PC is provided for controlling the lighting units through the addresses of the lighting units, stored in a signal control unit.\nThe recent development of controllable light sources increases the possibilities for a lighting designer. Due to this, the design of lighting scenes rises in significance by allowing to apply various effects and atmospheres without a change in the lighting units. Consequently, such a lighting design can be regarded as a complex work product and thus intellectual property of the designer.\nA lighting design is usually implemented in programs or scripts for operating the respective lighting units. An undesirable side-effect of such programs is that copying is generally possible without much effort, enabling use of such lighting design without consent of the designer.\nAccordingly, it is an object of this invention to provide a lighting system and a method for operating a lighting system which allows an efficient commercial use of lighting designs."}
{"text": "It is often desirable to be able to inject a liquid having certain characteristics into a flow of a second liquid having different characteristics. For example, in spraying water or fertilizer on a lawn it may be desirable to inject into the flow of water or fertilizer a second liquid, such as a herbicide or pesticide. Similarly, it may be desirable to inject a liquid detergent into a flow of water during a washing process.\nWhere the injection under pressure of one liquid into a flow of second liquid is desired, a separate pressurizing source for the injected liquid will usually be required."}
{"text": "Single server high availability (SSHA) storage systems have multiple controllers present in a topology, typically employing Redundant Array of Independent Disk methodologies. In some persistent reservation implementations, one storage controller has exclusive access to a particular set of logical volumes configured from a plurality of disk drives or other storage devices. Another storage controller in this implementation may have exclusive access to another particular set of logical volumes and be able to “see” logical volumes of its counterpart storage controller without being able to access those logical volumes, and vice versa. Because of this, a host driver of the server processing input/output (I/O) requests to the logical volumes may misdirect a request to a storage controller. Serial Attached Small Computer System Interface (SAS) storage controllers overcome this problem by rerouting or “shipping”, the I/O request and its associated data over a SAS connection between the storage controllers. But, this implementation requires much computational processing and time to perform the transfer. Cache mirroring (e.g., in a multipath I/O system with storage controller redundancy) is another computationally intensive and lengthy implementation where data is pulled by one storage controller responding to the I/O request and transferred to the other storage controller through the SAS connection between controllers."}
{"text": "1. Field of the Invention\nThis invention relates to system hardware detection and, more particularly, to the detection of system components and their features and capabilities.\n2. Description of the Related Art\nMany computer systems include one or more system management subsystems that may employ components which are coupled to the system management subsystem via some communication path or management bus such as an I2C bus for example. The components may be general-purpose devices such as general-purpose input/output (GPIO) devices or they may be task specific devices such as temperature measurement devices. In either case, the components may be simple devices with little or no processing power. In addition, the system management subsystem may include a controller that may identify each system management device and may also manage any information retrieved from such a system management device. Further, each component may reside at a fixed location or address on the management bus. Upon system initialization, the controller may access a component and feature set map which is hard-coded into the controller. The component and feature set map may identify each of the system management devices coupled to the management bus.\nHowever since the component and feature set map is hard-coded into a given controller, it may be difficult to interchange controllers from one system to another. In addition, it is possible that a component may be installed that either includes features not in a particular controller's feature set map or the component type is not in the controller's component map. In either case, the controller may not be able to access or control all of the component's features or may not be capable of communicating with that particular component at all."}
{"text": "Paper is widely used to display and record information. Printed information is easier to read than information displayed on a computer screen. Hand-drawing and handwriting afford greater richness of expression than input via a computer keyboard and mouse. Moreover, paper doesn't run on batteries, can be read in bright light, more robustly accepts coffee spills, and is portable and disposable.\nOnline publication has many advantages over traditional paper-based publication. From a consumer's point of view, information is available on demand, information can be navigated via hypertext links, searched and automatically personalized.\nFrom the publisher's point of view, the costs of printing and physical distribution are eliminated, and the publication becomes more attractive to the advertisers who pay for it because it can be targeted to specific demographics and linked to product sites.\nOnline publication also has disadvantages. Computer screens are inferior to paper. At the same quality as a magazine page, an SVGA computer screen displays only about a fifth as much information. Both CRTs and LCDs have brightness and contrast problems, particularly when ambient light is strong, while ink on paper, being reflective rather than emissive, is both bright and sharp in ambient light."}
{"text": "(a) Field of the Invention\nThe present invention relates to a display device and method thereof. More particularly, the present invention relates to a display device having improved display quality and a method of preventing signal delay in the display device.\n(b) Description of the Related Art\nA liquid crystal display (“LCD”), a plasma display panel (“PDP”), and an organic light emitting device (“OLED”) are among widely used flat panel displays.\nAmong these flat panel displays, the LCD and the OLED include a display panel including switching elements and display signal lines, and a gate driver outputting gate signals to gate lines among the display signal lines to turn on/off the switching elements. The gate driver may be made of at least one chip or may be integrated with the display panel.\nA signal controller that processes externally input image data and controls the gate driver uses a plurality of control signals for the control of the gate driver. The signal controller generates various signals such as a scanning start signal to instruct the application of the gate signal, a gate clock signal to determine the application time of the pulse width of the gate signal, and an output enable signal to determine a pulse width of the gate signal.\nWhen integrating the gate driver with the display panel, the signals for driving the gate driver are transmitted through the signal lines connected to a signal generator, and the gate driver transmits the gate signals to the gate lines."}
{"text": "A database-management means may comprise an optimizer engine that identifies an efficient way, known as an “execution plan,” to respond to a database-access request. Such an optimizer engine may be necessary when a database-management means can derive an execution plan from a large number of possible data-access strategies, data join methods, and other implementation-specific details when in order to serve a query, generate a report, or respond to a request to update stored data. In all but the most modest configurations, selecting the most efficient execution plan may require an optimizer engine to consider a huge number of possibilities.\nMost optimizers address such tasks by analyzing a set of statistics that characterize the contents of a database. These statistics may include information about stored database objects like tables, fields, records, indices or keys, and partitions. Collecting or gathering such statistics may be a resource-intensive task that requires a distinct statistics-collection function, and this function may need to be performed frequently in order to ensure that the statistics it collects are as current as possible or, as described herein, to ensure that the statistics have not become “stale.”\nEach such statistics-collection or statistics-generation job may consume so many resources that it reduces database performance to unacceptable levels. When possible, these functions may be scheduled as batch jobs that run after hours in order to minimize their effect on users.\nThese conflicting requirements may create a paradox for database system designers. If a mission-critical database must be available at all times, it may be difficult to identify idle periods in which a statistics-collection task may be run with minimal effect on database availability or performance. But if statistics-collection tasks are not performed frequently enough, the statistics available to an optimizer engine may become too stale to allow the optimizer to identify execution plans that consistently ensure a efficient response to a database-access request."}
{"text": "This invention relates to a method and system for increasing accuracy of clamping force of electric brakes of aircraft, and more particularly relates to a method and system for increasing accuracy of clamping force of electric aircraft carbon brakes providing greater accuracy for low brake clamping force commands by dedicating a portion of a plurality of electric brake actuators of each brake to low brake clamping force commands, without otherwise affecting normal braking. The present invention also relates to a system and method for aircraft brake metering to alleviate structural loading, and more particularly relates to a system and method for metering aircraft brakes to alleviate structural loading of an aircraft by delaying a full onset of braking for a preset period of time, such as on any brake-by-wire aircraft where a brake metering function can be modified.\nCommercial aircraft commonly have landing gear with electrically actuated brakes for wheels mounted to the wing and body of the aircraft. The electrically actuated brakes are typically carbon brakes including a torque plate and a carbon heat sink stack containing the friction surfaces that are clamped together by four electric brake actuators with a clamping brake force to cause a wheel to decrease its speed of rotation. In such a conventional airplane carbon brake system, when braking is commanded, either by a pilot's actuation of a brake pedal or automatic braking, it causes the friction surfaces of the carbon brakes to make contact, creating brake torque to slow down the rotational speed of the wheel, and through contact with the ground, the taxi speed of airplane.\nAs is described in U.S. Pat. No. 7,441,844, it is possible to reduce brake wear of electrically operated aircraft carbon brakes, once braking has been commenced, by maintaining a minimum light residual clamping brake force when braking is no longer commanded, such as when a pilot stops pressing on a brake pedal, or otherwise during a commanded release of braking during automatic braking. During taxiing of commercial aircraft, particularly at low speeds, steering of the aircraft is typically controlled by braking, and an unequal distribution of brake energy due to inaccurate metering of brake clamping force can in some instances interfere with the directional stability of aircraft, particularly when a minimum light residual clamping brake force is maintained during taxiing when braking is no longer commanded. Unequal distribution of brake energy due to inaccurate metering of brake clamping force can also result in damage to wheels and brakes from exposure to excessively high temperatures. It has been found that it is not possible with currently available electrical braking systems to achieve a brake clamping force accuracy required by current industry standards for the Boeing 787, and to prevent unequal distribution of brake energy.\nRepeated rapid brake applications at low speed can also create fatigue loads on aircraft structural components. This is most apparent when the aircraft is taxiing at low speed and full rapid braking is commanded by a pilot by rapidly depressing the brake pedals. On some aircraft, such a condition will cause the nose gear of the aircraft (which does not typically have brakes) to bounce off the runway, because sudden braking loads are presented as a torque about the fuselage and reacted as a down force on the nose gear of the aircraft. Over the life of an aircraft, the design of the aircraft must take into account these structural loads so that the aircraft structure will not fail prematurely.\nPrevious known methods for alleviating airplane structural load during rapid low speed braking typically have involved the delayed application of some, but not all of the available brakes to reduce the aircraft structural loads and fatigue. Such selective delayed application of less than all available aircraft brakes, such as by applying only some of the available brakes while the remainder of the available brakes are not applied during an initial brake application during taxi conditions, can reduce aircraft structural loads.\nOne such conventional electric brake system for an aircraft is known that employs a brake control process to reduce high dynamic structural loading of the aircraft landing gear and lurching of the aircraft that can be caused by braking maneuvers. The system obtains current aircraft speed to determine that the aircraft is in taxi mode, current brake pedal deflection position, and deflection rate to determine whether to delay the onset of a desired braking condition.\nIt would be desirable to provide a system and method for alleviating aircraft structural loads during braking that does not require a logic condition of first determining that the aircraft is in taxi mode, and that does not require releasing and then reapplying some, but not all of the brakes during low speed braking. It would be desirable to provide a system and method for metering aircraft brakes to alleviate structural loading of an aircraft that does not rely upon the delayed application of brakes, but instead delays a full onset of braking for a preset period of time, to reduce aircraft structural loading, such as on any brake-by-wire aircraft where a brake metering function can be modified. The present invention meets these and other needs."}
{"text": "1) Field of the Invention\nThis invention relates generally to rocking reclining chairs, and more particularly to a rocking reclining chair having a lock for preventing the chair from rocking while in reclined and unreclined positions.\n2) Description of Related Art\nReclining chairs are well known in the art and typically comprise a wood or metal frame defining a seat bottom, a seat back for cushioning a user\"\"s back, two armrests, and an extendable footrest. One type of recliner is a xe2x80x9cthree-wayxe2x80x9d reclining chair. A seated user in a three-way reclining chair can adjust the backrest into an upright position, a partly reclined xe2x80x9cTVxe2x80x9d position, or a fully reclined position. In the upright position, the seat back is upright and the footrest is retained in the base of the chair to allow the user\"\"s feet to rest on the floor. In the TV position, the footrest is extended and the seat back is reclined at a relatively shallow angle from the upright position, which still allows the user to watch television. In the fully reclined position the seating surfaces of the chair are nearly horizontal to the floor. Changing position from the upright to TV and fully reclined positions typically requires extension of the footrest using a lever positioned on an outer side of the armrest and pushing backward on the seat back while gripping the armrests.\nRocking reclining chairs, or rocker recliners, are reclining chairs that rock when the recliner is in an upright position. The frame of a rocker recliner typically includes a stationary base frame and a rocking frame. The rocking frame is attached to the seating surfaces and is supported by a pair of rocking cams. The rocking cams rest and reciprocate on the base frame to produce the rocking motion. The rocking frame and the base frame are also connected via a spring box which maintains contact between the rocking cams and the base frame during rocking and limits the rocking motion to prevent instability of the recliner. The spring box also biases the rocking cams, and hence the seating surfaces, to a neutral upright position when the rocker recliner is not being urged forward or backward during rocking. Due to potential instability of the recliner in the TV and fully reclined positions, most rocker recliners include a mechanism for automatically inhibiting rocking motion when in the reclined positions. Rocking is typically inhibited by a set of stops that extend down from the upper part of the rocking frame and onto the base frame or rocking cams as the footrest is extended.\nAlso known in the art are rocker recliners that include a lockout mechanism for manually inhibiting rocking motion, even when the rocker recliner is in the upright position. Although rocking the rocker recliner with the seat back in an upright position does not make the recliner unstable, seated users oftentimes find it difficult to rise from the chair while it is rocking. This is especially true for elderly users who lack the leg strength to stand unassisted and must grip the armrests to assist in the rising motion. The lockout mechanism, which can be controllable by a lever that moves a lock cam into contact with the rocking cams or base frame, allows the user to prevent rocking motion.\nDespite the many improvements that have been made in rocker recliners, further improvement has been needed in the respect of simplifying a rocker recliner in terms of operation, structure, and expense. For example, the use of a manual lockout mechanism instead of an automatic locking mechanism complicates the operation of the rocker recliner by requiring the user to operate the lockout mechanism to prevent rocking while in the reclined positions. Also, the addition of a lockout mechanism also necessitates costly additional components and assembly time. Additionally, the design is complicated, increasing the likelihood of assembly errors and, possibly, wear or failure of the chair.\nThus, there exists a need for a rocker recliner that provides stable reclining as well as rocking and no-rocking modes in the upright position. While in the no-rocking mode in the upright position, the rocker recliner should be stable so that a seated user can use the rocker recliner for support while standing. Specifically, the armrests should remain stationary so that the seated user can grip and push on the armrests while standing. Additionally, the rocker recliner should be simple to use, and the structure should be simple for the sake of manufacturing ease and cost.\nThe present invention provides a self-locking rocker recliner that addresses these needs in the prior art. The rocker recliner includes a dual-action locking mechanism that adjusts between an unlocked position in which the rocker recliner can rock and a locked position in which rocking is inhibited. The dual-action locking mechanism is adjusted automatically when the rocker recliner is reclined, but can also be adjusted manually by a lockout lever that is rotated by a user. Thus, the rocker recliner provides an automatic no-rocking mode for stability when the rocker recliner is reclined and a manual no-rocking mode whenever the user desires additional stability, for example while standing from or sitting into the chair.\nIn one embodiment of the present invention, the rocker recliner includes a base frame and a rocking frame that is rockably connected to said base frame. The rocking frame includes two rocking cams that rest and reciprocate on the base frame. Rocking is also facilitated by a spring box that biases the rocking frame to a neutral position. The rocking frame supports seating surfaces of the rocker recliner, such as a seat back, a seat bottom, a footrest, and armrests. The rocking frame preferably includes an upper rocking frame and a lower rocking frame that are movably connected so that the rocker recliner can be adjusted between an upright position and at least one reclined position, such as a TV or fully reclined position. The rocker recliner can be reclined by rotating a recline lever on the side of the rocker recliner, and reclining the rocker recliner can also cause the extension of a footrest. A lockout lever, which is attached to the rocker recliner, is configured to be moved by a user between a locked position and an unlocked position. The lockout lever is linked to a dual-action locking mechanism rotatably connected to the rocker recliner such that rotation of the lockout lever causes a corresponding rotation of the dual-action locking mechanism between a locked position and an unlocked position. Rotation of the dual-action locking mechanism to the locked position prevents rocking of the rocker recliner. The dual-action locking mechanism is also configured to rotate to the locked position when the rocking frame is reclined to the reclined position. In one embodiment, an over-center mechanism biases the dual-action locking mechanism to the locked and unlocked positions. In another embodiment, the lockout lever is located between an armrest and a seat bottom of the rocker recliner. This location provides convenient access for the user while maintaining the pleasing appearance of the rocker recliner.\nIn another embodiment of the present invention, the rocker recliner includes a base frame, a lower rocking frame, at least one seating surface attached to the lower rocking frame, and an upper rocking frame. The lower rocking frame is rockably connected to the base frame by at least one rocking cam. The rocking cam is supported by the base frame and defines a curved surface that defines an arcuate rocking motion of the lower rocking frame relative to the base frame. The upper rocking frame is movably connected to the lower rocking frame by a recline linkage configured to adjust the upper rocking frame between an upright position and a reclined position relative to the lower rocking frame. The rocker recliner also includes a dual-action locking member, a lockout linkage, a recline linkage, and a stop blocker. The dual-action locking member is rotatably connected to the lower rocking frame such that the dual-action locking member can pivot about a pivot pin between a locked position and an unlocked position. The dual-action cam plate can be biased to the locked and unlocked positions by an over-center mechanism. The lockout linkage is connected to the dual-action locking member and configured to rotate the dual-action locking member between the locked and unlocked positions. A lockout lever can attach to the lockout linkage so that moving the lockout lever to a locked position causes the dual-action locking member to be moved to the locked position thereof regardless of position of the recline linkage. The recline lock link extends from the recline linkage and is connected to the dual-action locking member such that adjusting the recline linkage to the reclined position causes the recline lock link to rotate the dual-action locking member to the locked position. When the dual-action locking member is in the locked position, the stop blocker is configured to contact the base frame. The stop blocker can be connected to a blocker plate that defines a slot for retaining a stud attached to the base frame. The dual-action locking member can include a dual-action cam plate that defines a cam plate slot that retains a stud attached to the recline lock link. Alternatively, a dual-action linkage can connect the dual-action locking member to the stop blocker and dual-action lock link. A slot connection connects the dual-action linkage to the recline linkage. For example, the dual-action linkage can define a slot that retains a stud attached to the recline lock link.\nThus, the present invention provides a rocker recliner which inhibits rocking while the rocker recliner is in a reclined position and, alternately, when a lockout lever is in a locked position. The lockout lever can be locked while the rocker recliner is in the upright position so that the rocker recliner provides stability to a user that is standing from or sitting into the rocker recliner. Specifically, the seating surfaces, such as the armrests, are held stationary so that the seated user can grip and push on the armrests while standing. Additionally, both reclining and the manual lockout cause rocking to be inhibited by a dual-action locking mechanism, thus simplifying and economizing the design and manufacture of the rocker recliner."}
{"text": "In recent years, in a field of materials processing with high-power lasers, attention is paid to processing by a laser apparatus using an optical fiber to which rare earth elements are doped. The laser emits a wavelength of 1 μm band, and an optical-fiber delivery is possible unlike CO2 laser used in the related art. Additionally, since this laser can be focused to a smaller spot than CO2 laser, it is suitable for high-speed processing or fine micro processing, such as cutting, welding, or the like. Regarding this laser apparatus, there are already a number of well-known techniques (for example, refer to Patent Document 1).\nAn example of a laser materials processing apparatus 100 of the related art is schematically shown in FIG. 13. The laser materials processing apparatus 100 is generally constituted by a laser apparatus 101, an optical fiber 122, and an irradiation optical system (a collimating lens 105 and a condenser lens 106). A laser beam 104, which is generated in the laser apparatus 101, is guided and transmitted through the optical fiber 122, and exits the emitting end 103, passes through the collimating lens 105 and the condenser lens 106, and reaches an irradiation point α on a workpiece 107.\nA related-art example of the above laser apparatus 101 is shown in FIG. 14. The laser apparatus 101 is generally constituted by an active fiber 120 that is a double-clad fiber in which rare earth elements are doped to a parent material of a core, fiber Bragg gratings (hereinafter abbreviated as FBGs) 121 that are formed near both ends of this active fiber 120 to act as mirrors of laser resonators, a multi-coupler 123, a plurality of semiconductor laser beam sources 126, and an optical fiber 122 that is connected to the active fiber 120 via a connecting point 129.\nThe double-clad fiber is an optical fiber having double claddings around a core. In the related art, the double-clad fiber is used as a laser medium by doping rare earth elements or the like to the parent material of the core. Inner cladding has two kinds of actions. One action confines the light within the core. The other action confines pumping light that pumps the laser medium within the core. Outer cladding has an action that confines the pumping light.\nThe pumping light emitted from the plurality of semiconductor laser beam sources 126 is coupled into the inner cladding of the active fiber 120 via the optical fiber 125 and the multi-coupler 123. The laser beam generated in the core of the active fiber 120 is amplified while going back and forth between both the FBGs 121, and a portion thereof is taken out from one FBG 121 (FBG 121 on the right of the plane of the sheet). This laser beam is coupled into the optical fiber 122 located downstream of the connecting point 129, and is emitted to exterior space from an output end 103. In addition, a passive optical fiber for transmission with single-layer cladding is used as the optical fiber 122 located downstream of the connecting point 129. As described above, the laser apparatus of the configuration using an active fiber as the laser medium is referred to as a fiber laser.\nThe power density distribution at a cross-section of the optical fiber 122 may be a single mode close to a Gaussian shape or a multi-mode close to a top-hat shape.\nA laser beam of an output of 2 kW or more is used for macro processing referred to as cutting of steel of a thickness of 1 mm or more. The single mode output in these region causes, the large optical loss due to nonlinear effects such as stimulated Brillouin scattering or Raman scattering within the core of the optical fiber. Therefore, the distance by which fiber transmission is allowed is limited to about several meters. Thus, in the laser materials processing apparatus, generally, a fiber laser of a multi-mode where a larger fiber transmission distance is secured is used.\nThe power density distribution of the laser beam at the irradiation point α is an image of the outlet (the emitting end 103 of the laser apparatus 101) of the optical fiber 122, formed by an irradiation optical system including the collimating lens 105 and the condenser lens 106. In the case of the fiber laser where the power density distribution in the optical fiber 122 being multi-mode, the power density distribution at the irradiation point α becomes a substantially uniform top-hat shape, as shown in FIG. 15, which is uniform and almost the same distribution within the optical fiber 122. In addition, in FIG. 15, X represents the distance from the center O of a laser beam at the irradiation point α, and I represents the power density of the laser beam."}
{"text": "It is well known that the quality of food is improved when it has been properly stirred or mixed during preparation. Stirring and mixing promotes even heating and cooking so that the finished product is consistently and thoroughly prepared. This is true whether the food is rice, oatmeal, meat, or popcorn. However, it is also known that the need to mix or stir food while it is being cooked is a time consuming, tiresome, and often tedious aspect of food preparation. Various attempts have been made to obviate or minimize the need for manual stirring of food while it is being cooked. These attempts include use a vibrating device to impart motion to the cooking system which causes relative movement between the system and the food being cooked as well as within the food itself.\nThe preparation of popcorn requires particular attention. In order to prepare a sufficient quantity of popped kernels that are uniformly cooked, the even application of heat to all of the kernels is imperative. Otherwise, some kernels pop earlier than others and become burned while the later-popping kernels continue to be heated. On the other hand, uneven heating can leave numerous kernels unpopped or only partially popped.\nMany previous attempts have been made to engineer an efficient way to prepare high quality popcorn. For example, the well known Jiffy Pop® brand combines popcorn kernels and oil in a disposable aluminum pan with an expandable aluminum foil top. The Jiffy Pop® unit is heated on an electric range while the consumer continuously oscillates the unit, causing the kernels to move and mix in the pan. One obvious disadvantage of this method is the nearly constant attention that the consumer must pay to the task of shaking the disposable pan across the heat source. The Jiffy Pop® container must be continuously shaken back and forth so that the kernels are continuously mixed while being heated.\nAnother example is the use of a microwave oven to heat a paper package containing kernels and oil. This method has well-known problems, including wildly uneven heating of the kernels resulting in burned popcorn and many kernels remaining unpopped. The overall flavor of popcorn prepared in the microwave is also well-known to be below average.\nThus, there is also a need for a reliable system to make high quality popcorn without requiring the entirety of the consumer's attention during preparation."}
{"text": "This application is a 371 of PCT/EP98/02878 filed May 15, 1998, now WO 98/54137 Dec. 3, 1998.\nThe present invention relates to novel substituted 2-phenylpyridines of the formula I \nin which the substituents and the index m have the following meanings:\nm is 0 or 1;\nis halogen, C1-C4-haloalkyl, C1-C4-haloalkoxy, C1-C4-alkylthio, C1-C4-alkylsulfinyl, C1-C4-alkylsulfonyl, C1-C4-haloalkylthio or cyano;\nR2 is fluorine or trifluoromethyl;\nR3 is hydrogen or halogen;\nR4 is halogen or cyano;\nR5 is CO2R6, OR7, SR7, C(R8)xe2x95x90Nxe2x80x94Oxe2x80x94R7 or C(R8)xe2x95x90C(R8)xe2x80x94COxe2x80x94Oxe2x80x94R6, where\nR6 is hydrogen, an unsubstituted or halogen-substituted C1-C8-alkyl-, C3-C6-alkenyl- or C3-C6-alkynyl radical; C1-C4-alkoxy-C1-C4-alkyl, C1-C6-alkoxycarbonyl-C1-C4-alkyl, C3-C4-alkenyloxycarbonyl-C1-C4-alkyl, C3-C4-alkynyloxycarbonyl-C1-C4-alkyl or C1-C4-alkoxy-(C1-C4-alkoxy)carbonyl-C1-C4-alkyl;\nR7 may have the meaning of R6 or may be CH2xe2x80x94CO2[C1-C4-alkylene)]xe2x80x94CO2R9 or CH[C1-C4-alkyl]CO2xe2x80x94[C1-C4-alkylene]xe2x80x94CO2R9;\nR8 is hydrogen, halogen or C1-C4-alkyl and\nR9 is hydrogen or C1-C4-alkyl,\nand the agriculturally useful salts of the compounds I.\nFurthermore, the invention relates to\na process for preparing the compounds I and\nintermediates of the formula II,\nherbicides and compositions for the desiccation and/or defoliation of plants which comprise the compounds I as active substances,\nmethods for controlling undesirable vegetation and for the desiccation and/or defoliation of plants using the compounds I.\nSubstituted 2-phenyl-3-chloropyridines having herbicidal activity are already known from WO 95/02580, WO 95/02590 and WO 97/11059.\nHowever, the herbicidal activity of the prior art compounds with respect to harmful plants is not always entirely satisfactory.\nIt is an object of the present invention to provide novel herbicidally active compounds which allow better selective control of undesirable plants. It is a further object to provide novel compounds which have desiccant/defoliant action.\nWe have found that these objects are achieved by the substituted 2-phenylpyridines of the formula I defined at the outset having herbicidal activity, and by the novel intermediates II for their preparation.\nDepending on the substitution pattern, the compounds of the formula I can contain one or more chiral centers, in which case they exist in the form of enantiomer or diastereomer mixtures. The invention relates to the pure enantiomers or diastereomers and also to mixtures thereof.\nThe substituted 2-phenylpyridines I where R6, R7 and R9=hydrogen may be present in the form of their agriculturally useful salts, the kind of salts usually not being important. Suitable in general are the salts of those bases whose herbicidal activity is not impaired in comparison to the free compound I.\nSuitable salts are in particular those of the alkali metals, preferably sodium salts and potassium salts, of the alkaline earth metals, preferably calcium salts and magnesium salts, those of the transition metals, preferably zinc salts and iron salts, and ammonium salts where the ammonium ion, if desired, may carry one to four C1-C4-alkyl or hydroxy-C1-C4-alkyl substituents and/or one phenyl or benzyl substituent, preferably diisopropylammonium, tetramethylammonium, tetrabutylammonium, trimethylbenzylammonium and trimethyl-(2-hydroxyethyl)ammonium salts, furthermore phosphonium salts, sulfonium salts such as preferably tri-(C1-C4-alkyl)sulfonium salts and sulfoxonium salts such as preferably tri-(C1-C4-alkyl)sulfoxonium salts.\nThe terms alkyl, alkylene haloalkyl, alkoxy, carboxyalkyl, alkoxyalkyl, alkoxycarbonylalkyl, alkenyl and alkynyl used in the definition of the substituents R1, R6, R7, R8 and R9 arexe2x80x94like the term halogenxe2x80x94collective terms for individual enumerations of the individual group members. All alkyl moieties may be straight-chain or branched. The haloalkyl radical preferably carries one to five identical or different halogen atoms.\nSpecific meanings are, for example:\nhalogen: fluorine, chlorine, bromine and iodine, preferably fluorine and chlorine;\nC1-C4-alkyl: methyl, ethyl, n-propyl, 1-methylethyl, n-butyl, 1-methylpropyl, 2-methylpropyl and 1,1-dimethylethyl;\nC1-C6-alkyl: C1-C4-alkyl as mentioned above, and n-pentyl, 1-methylbutyl, 2-methylbutyl, 3-methylbutyl, 2,2-dimethylpropyl, 1-ethylpropyl, n-hexyl, 1,1-dimethylpropyl, 1,2-dimethylpropyl, 1-methylpentyl, 2-methylpentyl, 3-methylpentyl, 4-methylpentyl, 1,1-dimethylbutyl, 1,2-dimethylbutyl, 1,3-dimethylbutyl, 2,2-dimethylbutyl, 2,3-dimethylbutyl, 3,3-dimethylbutyl, 1-ethylbutyl, 2-ethylbutyl, 1,1,2-trimethylpropyl, 1,2,2-trimethylpropyl, 1-ethyl-1-methylpropyl and 1-ethyl-2-methylpropyl;\nC1-C8-alkyl: C1-C6-alkyl as mentioned above, and, inter alia, n-heptyl, n-octyl;\nC3-C4-alkenyl: prop-1-en-1-yl, prop-2-en-1-yl, 1-methylethenyl, n-buten-1-yl, n-buten-2-yl, n-buten-3-yl, 1-methylprop-1-en-1-yl, 2-methylprop-1-en-1-yl, 1-methylprop-2-en-1-yl and 2-methylprop-2-en-1-yl;\nC3-C6-alkenyl: C3-C4-alkenyl as mentioned above, n-penten-1-yl, n-penten-2-yl, n-penten-3-yl, n-penten-4-yl, 1-methylbut-1-en-1-yl, 2-methylbut-1-en-1-yl, 3-methylbut-1-en-1-yl, 1-methylbut-2-en-1-yl, 2-methylbut-2-en-1-yl, 3-methylbut-2-en-1-yl, 1-methylbut-3-en-1-yl, 2-methylbut-3-en-1-yl, 3-methylbut-3-en-1-yl, 1,1-dimethylprop-2-en-1-yl, 1,2-dimethylprop-1-en-1-yl, 1,2-dimethylprop-2-en-1-yl, 1-ethylprop-1-en-2-yl, 1-ethylprop-2-en-1-yl, n-hex-1-en-1-yl, n-hex-2-en-1-yl, n-hex-3-en-1-yl, n-hex-4-en-1-yl, n-hex-5-en-1-yl, 1-methylpent-1-en-1-yl, 2-methylpent-1-en-1-yl, 3-methylpent-1-en-1-yl, 4-methylpent-1-en-1-yl, 1-methylpent-2-en-1-yl, 2-methylpent-2-en-1-yl, 3-methylpent-2-en-1-yl, 4-methylpent-2-en-1-yl, 1-methylpent-3-en-1-yl, 2-methylpent-3-en-1-yl, 3-methylpent-3-en-1-yl, 4-methylpent-3-en-1-yl, 1-methylpent-4-en-1-yl, 2-methylpent-4-en-1-yl, 3-methylpent-4-en-1-yl, 4-methylpent-4-en-1-yl, 1,1-dimethylbut-2-en-1-yl, 1,1-dimethylbut-3-en-1-yl, 1,2-dimethylbut-1-en-1-yl, 1,2-dimethylbut-2-en-1-yl, 1,2-dimethylbut-3-en-1-yl, 1,3-dimethylbut-1-en-1-yl, 1,3-dimethylbut-2-en-1-yl, 1,3-dimethylbut-3-en-1-yl, 2,2-dimethylbut-3-en-1-yl, 2,3-dimethylbut-1-en-1-yl, 2,3-dimethylbut-2-en-1-yl, 2,3-dimethylbut-3-en-1-yl, 3,3-dimethylbut-1-en-1-yl, 3,3-dimethylbut-2-en-1-yl, 1-ethylbut-1-en-1-yl, 1-ethylbut-2-en-1-yl, 1-ethylbut-3-en-1-yl, 2-ethylbut-1-en-1-yl, 2-ethylbut-2-en-1-yl, 2-ethylbut-3-en-1-yl, 1,1,2-trimethylprop-2-en-1-yl, 1-ethyl-1-methylprop-2-en-1-yl, 1-ethyl-2-methylprop-1-en-1-yl and 1-ethyl-2-methylprop-2-en-1-yl, preferably ethenyl and prop-2-en-1-yl;\nC3-C4-alkynyl: prop-1-in-1-yl, prop-2-in-3-yl, n-but-1-in-1-yl, n-but-1-in-4-yl, n-but-2-in-1-yl;\nC3-C6-alkynyl: C3-C4-alkynyl as mentioned above, and n-pent-1-in-1-yl, n-pent-1-in-3-yl, n-pent-1-in-4-yl, n-pent-1-in-5-yl, n-pent-2-in-1-yl, n-pent-2-in-4-yl, n-pent-2-in-5-yl, 3-methylbut-1-in-1-yl, 3-methylbut-1-in-3-yl, 3-methylbut-1-in-4-yl, n-hex-1-in-1-yl, n-hex-1-in-3-yl, n-hex-1-in-4-yl, n-hex-1-in-5-yl, n-hex-1-in-6-yl, n-hex-2-in-1-yl, n-hex-2-in-4-yl, n-hex-2-in-5-yl, n-hex-2-in-6-yl, n-hex-3-in-1-yl, n-hex-3-in-2-yl, 3-methylpent-1-in-1-yl, 3-methylpent-1-in-3-yl, 3-methylpent-1-in-4-yl, 3-methylpent-1-in-5-yl, 4-methylpent-1-in-1-yl, 4-methylpent-2-in-4-yl and 4-methylpent-2-in-5-yl, preferably prop-2-in-1-yl, 1-methylprop-2-in-1-yl;\nC1-C3-fluoroalkyl: C1-C3-alkyl as mentioned above where in each case 1-5 hydrogen atoms are replaced by fluorine, e.g. fluoromethyl, difluoromethyl, trifluoromethyl, 1-fluoroethyl, 2-fluoroethyl, 2,2-difluoroethyl, 2,2,2-trifluoroethyl, pentafluoroethyl, 3,3,3-trifluoropropyl, preference is given to difluoromethyl, trifluoromethyl, 2,2,2-trifluoroethyl, pentafluoroethyl and 3,3,3-trifluoropropyl, particular preference is given to trifluoromethyl;\nC1-C4-haloalkyl: C1-C4-alkyl as mentioned above which is partially or fully substituted by fluorine, chlorine and/or bromine, i.e. for example chloromethyl, dichloromethyl, trichloromethyl, fluoromethyl, difluoromethyl, trifluoromethyl, chlorofluoromethyl, dichlorofluoromethyl, chlorodifluoromethyl, 1-fluoroethyl, 2-fluoroethyl, 2,2-difluoroethyl, 2,2,2-trifluoroethyl, 2-chloro-2-fluoroethyl, 2-chloro-2,2-difluoroethyl, 2,2-dichloro-2-fluoroethyl, 2,2,2-trichloroethyl, pentafluoroethyl and 3-chloropropyl, preferably trifluoromethyl;\nC1-C4-haloalkoxy: C1-C4-alkoxy as mentioned above which is partially or fully substituted by fluorine, chlorine and/or bromine, i.e. for example chloromethoxy, dichloromethoxy, trichloromethoxy, fluoromethoxy, difluoromethoxy, trifluoromethoxy, chlorofluoromethoxy, dichlorofluoromethoxy, chlorodifluoromethoxy, 1-fluoroethoxy, 2-fluoroethoxy, 2,2-difluoroethoxy, 2,2,2-trifluoroethoxy, 2-chloro-2-fluoroethoxy, 2-chloro-2,2-difluoroethoxy, 2,2-dichloro-2-fluoroethoxy, 2,2,2-trichloroethoxy and pentafluoroethoxy, preferably C1-C2-haloalkoxy such as trifluoromethoxy;\nC1-C4-alkylthio: methylthio, ethylthio, n-propylthio, 1-methylethylthio, n-butylthio, 1-methylpropylthio, 2-methylpropylthio and 1,1-dimethylethylthio, preferably methylthio, ethylthio, methylethylthio;\nC1-C4-haloalkylthio: chloromethylthio, dichloromethylthio, trichloromethylthio, fluoromethylthio, difluoromethylthio, trifluoromethylthio, chlorofluoromethylthio, dichlorofluoromethylthio, chlorodifluoromethylthio, 1-fluoroethylthio, 2-fluoroethylthio, 2,2-difluoroethylthio, 2,2,2-trifluoroethylthio, 2-chloro-2-fluoroethylthio, 2-chloro-2,2-difluoroethylthio, 2,2-dichloro-2-fluoroethylthio, 2,2,2-trichloroethylthio and pentafluoroethylthio, preferably C1-C2-haloalkylthio such as trifluoromethylthio;\nC1-C4-alkylsulfonyl: methylsulfonyl, ethylsulfonyl, n-propylsulfonyl, 1-methylethylsulfonyl, n-butylsulfonyl, 1-methylpropylsulfonyl, 2-methylpropylsulfonyl and 1,1-dimethylethylsulfonyl;\nC1-C4-alkylsulfinyl: methylsulfinyl, ethylsulfinyl, n-propylsulfinyl, 1-methylethylsulfinyl, n-butylsulfinyl, 1-methylpropylsulfinyl, 2-methylpropylsulfinyl and 1,1-dimethylethylsulfinyl;\nC1-C4-alkoxy-C1-C4-alkyl: C1-C4-alkyl which is substituted by C1-C4-alkoxy such as methoxy, ethoxy, n-propoxy, 1-methylethoxy, n-butoxy, 1-methylpropoxy, 2-methylpropoxy and 1,1-dimethylethoxy, i.e. for example CH2OCH3, CH2OC2H5, n-propoxymethyl, (1-methylethoxy)methyl, n-butoxymethyl, (1-methylpropoxy)methyl, (2-methylpropoxy)methyl, (1,1-dimethylethoxy)methyl, 2-(methoxy)ethyl, 2-(ethoxy)ethyl, 2-(n-propoxy)ethyl, 2-(1-methylethoxy)ethyl, 2-(n-butoxy)ethyl, 2-(1-methylpropoxy)ethyl, 2-(2-methylpropoxy)ethyl, 2-(1,1-dimethylethoxy)ethyl, 2-(methoxy)propyl, 2-(ethoxy)propyl, 2-(n-propoxy)propyl, 2-(1-methylethoxy)propyl, 2-(n-butoxy)propyl, 2-(1-methylpropoxy)propyl, 2-(2-methylpropoxy)propyl, 2-(1,1-dimethylethoxy)propyl, 3-(methoxy)propyl, 3-(ethoxy)propyl, 3-(n-propoxy)propyl, 3-(1-methylethoxy)propyl, 3-(n-butoxy)propyl, 3-(1-methylpropoxy)propyl, 3-(2-methylpropoxy)propyl, 3-(1,1-dimethylethoxy)propyl, 2-(methoxy)butyl, 2-(ethoxy)butyl, 2-(n-propoxy)butyl, 2-(1-methylethoxy)butyl, 2-(n-butoxy)butyl, 2-(1-methylpropoxy)butyl, 2-(2-methylpropoxy)butyl, 2-(1,1-dimethylethoxy)butyl, 3-(methoxy)butyl, 3-(ethoxy)butyl, 3-(n-propoxy)butyl, 3-(1-methylethoxy)butyl, 3-(n-butoxy)butyl, 3-(1-methylpropoxy)butyl, 3-(2-methylpropoxy)butyl, 3-(1,1-dimethylethoxy)butyl, 4-(methoxy)butyl, 4-(ethoxy)butyl, 4-(n-propoxy)butyl, 4-(1-methylethoxy)butyl, 4-(n-butoxy)butyl, 4-(1-methylpropoxy)butyl, 4-(2-methylpropoxy)butyl or 4-(1,1-dimethylethoxy)butyl, preferably n-propoxymethyl, (1-methylethoxy)methyl, 2-(n-propoxy)ethyl and 2-(1-methylethoxy)ethyl and particularly preferably CH2OCH3, CH2OC2H5, 2-methoxyethyl or 2-ethoxyethyl;\n(C1-C6-alkoxy)carbonyl-C1-C2-alkyl: C1-C4-alkyl which is substituted by (C1-C6-alkoxy)carbonyl such as COOCH3, COOC2H5, n-propoxycarbonyl, COOCH(CH3)2, n-butoxycarbonyl, 1-methylpropoxycarbonyl, 2-methylpropoxycarbonyl, COOC(CH3)3, n-pentoxycarbonyl, 1-methylbutoxycarbonyl and n-hexoxycarbonyl, i.e. for example CH2xe2x80x94COOCH3, CH2xe2x80x94COOC2H5, n-propoxycarbonylmethyl, CH2xe2x80x94COOCH(CH3)2, n-butoxycarbonylmethyl, (1-methylpropoxycarbonyl)methyl, (2-methylpropoxycarbonyl)methyl, CH2xe2x80x94COOC(CH3)3, n-pentoxycarbonylmethyl, (1-methylbutoxycarbonyl)-methyl, n-Hexoxycarbonylmethyl, 1-(methoxycarbonyl)ethyl, 1-(ethoxycarbonyl)ethyl, 1-(n-propoxycarbonyl)ethyl, 1-(1-methylethoxycarbonyl)ethyl, 1-(n-butoxycarbonyl)ethyl, 1-(n-pentoxycarbonyl)ethyl, 1-(1-methylbutoxycarbonyl)ethyl, 1-(n-hexoxycarbonyl)ethyl, 2-(methoxycarbonyl)ethyl, 1-(n-hexoxycarbonyl)ethyl, 2-(methoxycarbonyl)ethyl, 2-(ethoxycarbonyl)ethyl, 2-(n-propoxycarbonyl)ethyl, 2-(1-methylethoxycarbonyl)ethyl, 2-(n-butoxycarbonyl)ethyl, 2-(1-methylpropoxycarbonyl)ethyl, 2-(2-methylpropoxycarbonyl)ethyl, 2-(1,1-dimethylethoxycarbonyl)ethyl;\n(C1-C6-alkoxy)carbonyl-C1-C4-alkyl: (C1-C6-alkoxy)carbonyl-C1-C2-alkyl as mentioned above, and 2-(methoxycarbonyl)propyl, 2-(ethoxycarbonyl)propyl, 2-(n-propoxycarbonyl)propyl, 2-(1-methylethoxycarbonyl)propyl, 2-(n-butoxycarbonyl)propyl, 2-(1-methylpropoxycarbonyl)propyl, 2-(2-methylpropoxycarbonyl)propyl, 2-(1,1-dimethylethoxycarbonyl)propyl, 3-(methoxycarbonyl)propyl, 3-(ethoxycarbonyl)propyl, 3-(n-propoxycarbonyl)propyl, 3-(1-methylethoxycarbonyl)propyl, 3-(n-butoxycarbonyl)propyl, 3-(1-methylpropoxycarbonyl)propyl, 3-(2-methylpropoxycarbonyl)propyl, 3-(1,1-dimethylethoxycarbonyl)propyl, 2-(methoxycarbonyl)butyl, 2-(ethoxycarbonyl)butyl, 2-(n-propoxycarbonyl)butyl, 2-(1-methylethoxycarbonyl)butyl, 2-(n-butoxycarbonyl)butyl, 2-(1-methylpropoxycarbonyl)butyl, 2-(2-methylpropoxycarbonyl)butyl, 2-(1,1-dimethylethoxycarbonyl)butyl, 3-(methoxycarbonyl)butyl, 3-(ethoxycarbonyl)butyl, 3-(n-propoxycarbonyl)butyl, 3-(1-methylethoxycarbonyl)butyl, 3-(n-butoxycarbonyl)butyl, 3-(1-methylpropoxycarbonyl)butyl, 3-(2-methylpropoxycarbonyl)butyl, 3-(1,1-dimethylethoxycarbonyl)butyl, 4-(methoxycarbonyl)butyl, 4-(ethoxycarbonyl)butyl, 4-(n-propoxycarbonyl)butyl, 4-(1-methylethoxycarbonyl)butyl, 4-(n-butoxycarbonyl)butyl, 4-(1-methylpropoxycarbonyl)butyl, 4-(2-methylpropoxycarbonyl)butyl or 4-(1,1-dimethylethoxycarbonyl)butyl, preferably CH2xe2x80x94COOCH3, CH2xe2x80x94COOC2H5, 1-(methoxycarbonyl)ethyl or 1-(ethoxycarbonyl)ethyl;\nC1-C3-alkoxy-(C1-C3-alkoxy)carbonyl-C1-C2-alkyl: (C1-C3-alkoxy)carbonyl-C1-C2-alkyl such as CH2COOCH3, CH2COOC2H5, CH2COOCH2xe2x80x94C2H5, CH2COOCH(CH3)2, CH(CH3)COOCH3, CH(CH3)COOC2H5, CH2CH2COOCH3, CH2CH2COOC2H5, CH2CH2COOCH2xe2x80x94C2H5, CH2CH2COOCH(CH3)2, 2-(COOCH3)propyl, 2-(COOC2H5)propyl, 2-(COOCH2xe2x80x94C2H5)propyl, 2-[COOCH(CH3)2]propyl, 3-(COOCH3)propyl, 3-(COOC2H5)propyl, 3-(COOCH2xe2x80x94C2H5)propyl, 3-[COOCH(CH3)2]propyl, preferably CH2COOCH3 or CH2COOC2H5, which is substituted in the C1-C3-alkoxy moiety by OCH3, OC2H5, OCH2xe2x80x94C2H5 or OCH(CH3)2, i.e. for example CH2COOCH2OCH3, CH2COOCH2OC2H5, CH2COOCH2OCH(CH3)2 or CH2COOCH2OC(CH3)3;\nC1-C4-alkoxy-(C1--C4-alkoxy)carbonyl-C1-C4-alkyl: C1-C3-alkoxy-(C1-C3-alkoxy)carbonyl-C1-C2-alkyl as mentioned above where one or both of the alkoxy moieties may additionally be n-butoxy, sec-butoxy, iso-butoxy or tert-butoxy, and 2-(COOCH3)butyl, 2-(COOC2H5)butyl, 2-(COOCH2xe2x80x94C2H5)butyl, 2-[COOCH(CH3)2]butyl, 3-(COOCH3)butyl, 3-(COOC2H5)butyl, 3-(COOCH2xe2x80x94C2H5)butyl, 3-[COOCH(CH3)2]butyl, 4-(COOCH3)butyl, 4-(COOC2H5)butyl, 4-(COOCH2xe2x80x94C2H5)butyl, 4-[COOCH(CH3)2]butyl.\nC1-C4-alkylene is, for example, methylene, 1,1-ethylene, 1,2-ethylene, 1,1-propylene, 1,2-propylene, 1,3-propylene, 2,2-propylene, 1,1-butylene, 1,2-butylene, 1,3-butylene, 1,4-butylene, 2,2-butylene, 2,3-butylene, 2-methyl-1,1-propylene, 2-methyl-1,2-propylene or 2-methyl-1,3-propylene, preferably methylene, 1,1-ethylene or 2,2-propylene.\nWith respect to the use of the substituted 2-phenylpyridines I according to the invention as herbicides and/or as desiccants/defoliants, the substituents and the index m preferably have the following meanings, in each case either alone or in combination:\nm is 0,\nR1 is C1-C3-fluoroalkyl, chlorine, methylsulfonyl or cyano;\nR2 is fluorine or trifluoromethyl;\nR3 is fluorine or chlorine;\nR4 is chlorine;\nR5 is CO2R6, OR7 or SR7, where\nR6 is hydrogen, C1-C5-alkyl, C3-C4-alkenyl, 3-chloroprop-2-ene, C3-C4-alkynyl, C1-C3-alkoxy-C1-C2-alkyl, C1-C6-alkoxycarbonyl-C1-C2-alkyl, propargyloxycarbonyl-C1-C2-alkyl, C1-C3-alkoxy-C1-C3-alkoxycarbonyl-C1-C2-alkyl;\nR7 may have the meaning of R6 or may be CH2xe2x80x94CO2[C1-C2-alkylene]xe2x80x94CO2R9 or CH[C1-C2-alkyl]xe2x80x94CO2xe2x80x94[C1-C2-alkylene]xe2x80x94CO2R9;\nR8 is hydrogen, halogen or C1-C4-alkyl and\nR9 is hydrogen or C1-C4-alkyl,\nand the agriculturally useful salts of the compounds I.\nParticular preference is given to the substituted 2-phenylpyridines Ia ({circumflex over (=)} where m=0, R2=fluorine and R4=chlorine), in particular to those compounds listed in Table A below:\nFurthermore, preference is given to the substituted 2-phenylpyridines Ib ({circumflex over (=)} I where m=0, R2=trifluoromethyl, R4=chlorine), in particular to the compounds listed in Table B below:\nThe substituted 2-phenylpyridines are obtainable by various routes, for example by the processes described in WO 95/02580 and WO 97/11059. The preparation of corresponding 2-phenylpyridine N-oxides of the formula I (m=1) can be carried out by the method of the process described in WO 97/11059. A more recent method of coupling the pyridine and phenyl components, according to which the pyridine sulfoxides IIb and pyridine sulfones IIc are reacted with Grignard or zinc reagents III or IV to give the final products I according to the invention, was described in DE Appl. No. 196 36995.9. \nIn the formulae III and IV, R3 to R5 are each as defined in claim 1 and X is in each case a halogen atom. The intermediates III and IV and the preparation thereof are described in DE Appl. No. 196 36995.9. \nThe thiopyridines II can be prepared by the method of the process described in DE Appl. No. 196 36997.5. DE Appl. No. 19722661.2 discloses a particularly favorable route to thiopyridines II starting from 2-halopyridines V and thio compounds of the formula VI in the presence of a copper catalyst.\nIn the formula II, R1 and R2 are each as defined in claim 1. The pyridine thioethers of the formula IIa (n=0) are starting materials for the preparation of the pyridine sulfoxides IIb (n=1) and pyridine sulfones IIc (n=2); the last two compounds being employed in the coupling reaction with III or IV. Z is a C1-10-alkyl, C2-C10-alkenyl or C2-C10-alkynyl radical with or without substitution by halogen, C1-C4-alkoxy, C1-C4-alkoxycarbonyl, di-(C1-C4-alkylamino)carbonyl, cyano or nitro, or is a C3-C8-cycloalkyl radical or a C1-C4-alkylenephenyl, phenyl or naphthyl radical with or without substitution in the phenyl moiety by halogen, C1-C3-alkyl, C1-C3-alkoxy, trifluoromethyl, cyano or nitro.\nThe terms alkyl, alkenyl, alkynyl, alkylene, alkoxy, alkoxycarbonyl, dialkylaminocarbonyl and cycloalkyl used in the definition of the substituent Z are collective terms for individual enumerations of the individual group members. All alkyl moieties may be straight-chain or branched. The haloalkyl radical preferably carries one to five identical or different halogen atoms.\nSpecific examples are:\nC1-C10-alkyl: C1-C8-alkyl as mentioned in the definition of substituents for R6, and n-nonyl and n-decyl;\n1-phenyl with or without substitution by halogen, C1-C3-alkyl, C1-C3-alkoxy, trifluoromethyl, cyano or nitro: 2-, 3-, 4-chlorophenyl, 2-, 3-, 4-tolyl, 2-chloro-4-methylphenyl, 2,4-dichlorophenyl, 2,4,6-trichlorophenyl, 2,6-dichloro-4-methylphenyl, 2-, 3-, 4-methoxyphenyl, 2-chloro-4-methoxyphenyl, 3-chloro-4-methoxyphenyl, 2-, 3-, 4-trifluoromethylphenyl, 2-, 3-, 4-cyanophenyl, 2-, 3-, 4-nitrophenyl, 2-methyl1-4-nitrophenyl, 2-chloro-4-trifluoromethylphenyl, 2-chloro-4-nitrophenyl and unsubstituted phenyl.\nParticular preference is given to those compounds II in which \nn is 1 or 2;\nR1 is trifluoromethyl, chlorine, methylsulfonyl or cyano;\nR2 is fluorine or trifluoromethyl and\nZ is an unsubstituted or chlorine- or methoxy-substituted C1-C8-alkyl radical, or a benzyl or phenyl radical which is unsubstituted or halogen-, methyl-, C1-C3-alkoxy-, trifluoromethyl-, cyano- or nitro-substituted in the phenyl moiety.\nSpecifically, mention may be made, for example, of the following pyridine thioethers IIa of Tables 1-4, of the pyridine sulfoxides IIb of Table 5-8 and of the pyridine sulfones IIc of Tables 9-12.\nPreference is given to the pyridine thioethers II.001-II.116 mentioned in Table 1 of Formula IIa1\nFurthermore, preference is given to the pyridine thioethers IIa2.001-IIa2.116 of the formula IIa2, which differ from the compounds IIa1.001-IIa1.116 in that in position 5 of the pyridine ring, a trifluoromethyl group replaces chlorine.\nFurthermore, preference is given to the pyridine thioethers IIa3.001-IIa3.116 of the formula IIa3, which differ from the compounds IIa1.001-IIa1.116 in that position 5 of the pyridine ring accommodates a methylsulfonyl group.\nFurthermore, preference is given to the pyridine thioethers IIa4.001-IIa4.116 of the formula IIa4, which differ from the compounds IIa1.001-IIa1.116 in that in position 5 of the pyridine ring a cyano group replaces chlorine.\nFurthermore, preference is given to the thiopyridines IIb1.001-IIb1.116 of the formula IIb1, which differ from the compounds IIa1.001-IIa1.116 in that they are the corresponding sulfoxides.\nFurthermore, preference is given to the thiopyridines IIb2.001-IIb2.116 of the formula IIb2, which differ from the compounds IIa2.001-IIa2.116 in that they are the corresponding sulfoxides.\nFurthermore, preference is given to the thiopyridines IIb3.001-IIb3.116 of the formula IIb3, which differ from the compounds IIa3.001-IIa3.116 in that they are the corresponding sulfoxides.\nFurthermore, preference is given to the thiopyridines IIb4.001-IIb4.116 of the formula IIb4, which differ from the compounds IIa4.001-IIa4.116 in that they are the corresponding sulfoxides.\nFurthermore, preference is given to the thiopyridines IIc1.001-IIc1.116 of the formula IIc1, which differ from the compounds IIa1.001-IIa1.116 in that they are the corresponding sulfones.\nFurthermore, preference is given to the thiopyridines IIc2.001-IIc2.116 of the formula IIc2, which differ from the compounds IIa2.001-IIa2.116 in that they are the corresponding sulfones.\nFurthermore, preference is given to the thiopyridines IIc3.001-IIc3.116 of the formula IIc3, which differ from the compounds IIa3.001-IIa3.116 in that they are the corresponding sulfones.\nFurthermore, preference is given to the thiopyridines IIc4.001-IIc4.116 of the formula IIc4, which differ from the compounds IIa4.001-IIa4.116 in that they are the corresponding sulfones.\nIf the compounds IIa-c are prepared using 2,3-difluoro-5-trifluoromethylpyridine and thiophenol as nucleophile and using hydrogen peroxide as oxidizing agent, the reaction can be illustrated by the following scheme: \nIt is also possible to use peracetic acid, sodium hypochlorite or chlorine and bromine as oxidizing agents instead of hydrogen peroxide in processes analogous to the above scheme.\nPreferred embodiments of the process are specified below.\nThe reaction of the 2-halopyridines V with a thiol VI is advantageously carried out in the presence of a solvent at 80-250xc2x0 C., preferably 120-200xc2x0 C., particularly preferably 140-180xc2x0 C.\nSolvents that are used for these reactionsxe2x80x94depending on the temperature rangexe2x80x94are hydrocarbons such as toluene and xylene, chlorinated hydrocarbons such as 1,2-dichloroethane, 1,1,2,2-tetrachloroethane, chlorobenzene, 1,2-, 1,3- or 1,4-dichlorobenzene, ethers such as 1,4-dioxane or anisol, glycol ethers such as dimethyl glycol ether, diethyl glycol ether or diethylene glycol dimethyl ether, esters such as ethyl acetate, propyl acetate, methyl isobutyrate or isobutyl acetate, carboxamides such as DMF or N-methylpyrrolidone, nitrated hydrocarbons such as nitrobenzene, ureas such as tetraethylurea, tetrabutylurea, dimethylethyleneurea and dimethylpropyleneurea, sulfoxides such as dimethyl sulfoxide, sulfones such as dimethyl sulfone, diethyl sulfone and tetramethylene sulfone, nitriles such as acetonitrile, propionitrile, butyronitrile or isobutyronitrile; water, or else mixtures of individual solvents.\nThe reaction is particularly preferably carried out in the melt without the use of a solvent.\nThe molar ratios in which the starting materials are reacted with each other are generally 0.9-1.4, preferably 0.95-1.1, particularly preferably 0.98-1.04, for the ratio of thiol to 2-halopyridine V. The concentration of the starting materials in the solvent is 0.1-5 mol/l, preferably 0.2-2 mol/l.\nThe reaction is promoted by the presence of a copper catalyst. Suitable catalysts are copper oxide, salts such as copper(II) chloride, copper sulfate, copper nitrate, copper acetate and copper carbonate. Particular preference is given to using finely distributed metallic copper, for example copper powder or copper bronze. The molar amount of catalyst, based on the 2-halopyridine V, is 0.001-10, preferably 0.001-1, particularly preferably 0.001-0.1 mol %.\nThe reaction can also be carried out in the presence of an organic base such as, for example, triethylamine, tri-n-propylamine, N-ethyldiisopropylamine, pyridine, xcex1-, xcex2-, xcex3-picoline, 2,4-, 2,6-lutidine, N-methylpyrrolidine, triethylene diamine, dimethylaniline, N,N-dimethylcyclohexylamine, quinoline or acridine.\nThe reaction is preferably carried out under acidic conditions by flushing the hydrogen halide that is eliminated during the reaction out of the reaction mixture by means of an inert gas, for example nitrogen, or by letting it escape into a gas washer under autogenous pressure.\nAdvantageously, the 2-halopyridine V is added over a period of 10 to 60 min to a mixture of the thiol VI and the catalyst at 20-80xc2x0 C., and the mixture is then stirred for another 0.5 to 12 hours, preferably 1 to 8 hours, at 140-180xc2x0 C. to allow the reaction to go to completion.\nHowever, it is also possible to add the thiol VI to a mixture of 2-halopyridine V and catalyst and then to complete the reaction as described above.\nIn the case of low-boiling 2-halopyridines V or thiols VI, the reaction can also be carried out in an autoclave.\nIf only one of the two starting materials has a low boiling point, the higher-boiling component can be initially charged together with the catalyst and the low-boiling component can be introduced directly at the reaction temperature of preferably 120-200xc2x0 C., particularly preferably 140-180xc2x0 C., or as a gas, at the rate of its consumption.\nThe reaction can be carried out under atmospheric or superatmospheric pressure, continuously or batchwise.\nThe oxidation of the pyridine thioethers of the formula IIa to the pyridine sulfoxides IIb and pyridine sulfones IIc is advantageously carried out with hydrogen peroxide, the pyridine sulfoxides IIb being obtained with approximately equivalent amounts of oxidant, and the pyridine sulfones IIc being obtained with approximately double the molar amount.\nExamples of solvents which can be used include water, acetonitrile, carboxylic acids such as acetic acid, trifluoroacetic acid, propionic acid, alcohols such as methanol, ethanol, isopropanol, tert-butanol and chlorinated hydrocarbons such as methyl ethyl ketone. Water, methanol, acetic acid and trifluoroacetic acid are particularly preferred.\nIn a particularly preferred variant, the reaction can also be catalyzed by adding relatively strong acids such as trifluoroacetic acid or perchloric acid. However, suitable catalysts additionally include metal compounds, eg. transition metal oxides such as vanadium pentoxide, sodium tungstate, potassium dichromate, iron oxide tungstate, sodium tungstate/molybdic acid, osmic acid, titanium trichloride, selenium dioxide, phenylselenenic acid, vanadyl 2,4-pentanedionate.\nThe catalysts are generally employed in an amount of from 0.5 to 10%, but it is also possible to employ stoichiometric amounts because the inorganic catalysts can easily be filtered off and recovered.\nAnother preferred oxidizing agent is peracetic acid or hydrogen peroxide/acetic anhydride, possibly also the peracetic acid which is present in equilibrium in a hydrogen peroxide/acetic acid mixture.\nAnother preferred oxidizing agent is pertrifluoroacetic acid or the hydrogen peroxide/trifluoroacetic acid mixture or else the hydrogen peroxide/trifluoroacetic anhydride mixture.\nOxidation with hydrogen peroxide in glacial acetic acid is generally very selective, but frequently slow. The reaction time can generally be reduced by adding trifluoroacetic acid. The oxidation with hydrogen peroxide in pure trifluoroacetic acid frequently leads to the formation of the corresponding N-oxides, as described for example in Chimia 29 (1975), 466. A rapid and selective oxidation of the pyridine thioethers IIa to the corresponding sulfoxides IIb and sulfones IIc is possible using solutions of hydrogen peroxide in mixtures of acetic acid and trifluoroacetic acid in the ratio of 10:1 to 1:1, in particular 6:1 to 4:1, by volume. Therefore, particular preference is given to using these mixtures as solvent.\nIt is possible furthermore to use as solvent petroleum ether, the abovementioned solvents and the abovementioned catalysts.\nBesides peracetic acid and pertrifluoroacetic acid, it is also possible to employ perbenzoic acid, monoperphthalic acid or 3-chloroperbenzoic acid, expediently in chlorinated hydrocarbons such as methylene chloride or 1,2-dichloroethane.\nAlso very suitable for oxidizing the thiols to sulfoxides or sulfones are chlorine and bromine. Favorable solvents are water, acetonitrile, dioxane, two-phase systems such as aqueous potassium bicarbonate solution/dichloromethane, and, in the case of pyridine alkyl thioethers, also acetic acid.\nIt is furthermore possible to employ as source of active halogen tert-butyl hypochlorite, hypochlorous and hypobromous acids, their salts, also N-halo compounds such as N-bromo- and N-chlorosuccinimide or else sulfuryl chloride.\nAlso favorable for the oxidation are dinitrogen tetroxide, eg. in the technically simple variant with air/nitrogen dioxide or trioxide and, for example, osmium(VIII) oxide as catalyst. The oxidation can also be carried out directly with nitric acid, in which case suitable additional solvents are acetic anhydride and acetic acid, and suitable catalysts are copper(I) and (II) bromide and chloride.\nAlso suitable for the oxidation is photosensitized oxygen transfer, in which case recommended photosensitizers are chlorophyll, protoporphyrin, Rose Bengal or Methylene Blue. Suitable inert solvents are hydrocarbons such as pentane, hexane, heptane, cyclohexane, chlorinated hydrocarbons such as methylene chloride, 1,2-dichloroethane, 1,1,2,2-tetrachloroethane, alcohols such as methanol, ethanol, n-propanol or isopropanol, ketones such as acetone, methyl ethyl ketone, polar aprotic solvents such as acetonitrile, propionitrile or aromatic hydrocarbons such as benzene, toluene, chlorobenzene or xylene. In place of oxygen, it is also possible to use ozone in the abovementioned solvents, plus ether, 1,4-dioxane or THF.\nBesides photosensitization, catalysts can also be recommended for oxidation with oxygen, eg. oxides and sulfides of nickel, copper, aluminum, tungsten, chromium, vanadium, ruthenium, titanium, manganese, molybdenum, magnesium and iron.\nEither pyridine sulfoxides IIb or their pyridine sulfones IIc are obtained depending on the stoichiometry of the oxidizing agents used. The molar ratios in which the starting materials are reacted with each other are generally 0.9-1.8, preferably 1.05-1.3, for the ratio of pyridine thioether IIa to oxidizing agent in the case of oxidation to pyridine sulfoxide IIb and generally 1.9-3.5, preferably 2.05-2.9, in the case of oxidation to pyridine sulfone IIC.\nThe concentration of the starting materials in the solvent is generally 0.1-5 mol/l, preferably 0.2-2 mol/l.\nIt is advantageous to introduce the pyridine thioether or the pyridine sulfoxide, if appropriate with one of the abovementioned catalysts, into one of the abovementioned solvents and then to add the oxidizing agent over the course of 0.25-20 hours with stirring. The addition and reaction temperatures depend on the optimum efficiency of the oxidizing agents in question and the suppression of side reactions. If photosensitized oxygen is used, the reaction is generally carried out at xe2x88x9220 to 80xc2x0 C., but in the case of metal catalysis, the reaction is generally carried out at 50 to 140xc2x0 C., and if ozone is used, the reaction is generally carried out at xe2x88x9278 to 60xc2x0 C. Owing to the limited solubility of the oxygen derivatives, they have to be introduced continuously over a prolonged period of time (up to 20 hours) into the reaction mixture until the oxidation has ended at the sulfoxide or sulfone stage. If air/nitrogen dioxide or trioxide are used, the reaction is preferably carried out at 15-150xc2x0 C. over a period of 1-15 hours. Liquid or easily soluble oxidizing agents such as hydrogen peroxide, the peracetic acid and pertrifluoroacetic acid which are formed together with acetic anhydride or in equilibrium with acetic acid and trifluoroacetic acid, respectively, hypochlorous or hypobromous acid, tert-butyl hypochlorite, chlorine or bromine, N-chlorosuccinimide or N-bromosuccinimide or nitric acid can be added over shorter periods of time of 0.25-6 hours to the reaction mixture of the pyridine thioether or sulfoxide, depending on the exothermic character of the reaction, to end the reaction after a further 1-60 hours. Moreover, preference is given to adding the liquid or dissolved oxidizing agent in portions. If hydrogen peroxide and peracetic acid or pertrifluoroacetic acid are used, the reaction is generally carried out at 0-90xc2x0 C., if tert-butyl hypochlorite is used, the reaction is generally carried out at xe2x88x9278 to 30xc2x0 C., if N-halogen compounds are used, the reaction is usually carried out at 0-30xc2x0 C., and if nitric acid is used, the reaction is usually carried out at 20 to 140xc2x0 C. If chlorine or bromine is used, a reaction temperature of 0-40xc2x0 C. is recommended.\nThe oxidations can be carried out under atmospheric or superatmospheric pressure, continuously or batchwise.\nAdvantageously, the multi-step reaction can also be carried out as a one-pot reaction, by reacting the thioethers IIa which are obtained in the first step of the synthesis in the reaction of the 2-halopyridines V with the thiols VI without isolation and purification directly to give the sulfoxides IIb or the sulfones IIc. Thus, if appropriate, the reaction product IIa is allowed to cool to 90-120xc2x0 C., a solvent, for example trifluoroacetic acid, preferably acetic acid and/or water, is added, if appropriate, and the oxidizing agent is then added at the rate of its consumption. Preferred oxidizing agents are hydrogen peroxide and especially sodium hypochlorite.\nFor work-up, the intermediates IIa-c are taken up in a water-imiscible solvent, acidic impurities and/or oxidizing agents are extracted using dilute bases or water, the solution is dried and the solvent is removed under reduced pressure.\nThe substituted 2-phenylpyridines I are usually preparable by one of the abovementioned processes. However, for economical or technical reasons it may be more advantageous to prepare some of the compounds I from similar 2-phenylpyridines which differ in the meaning of one radical.\nWork-up of the reaction mixtures is usually carried out by methods known per se, for example by diluting the reaction solution with water and subsequently isolating the product by filtration, crystallization or solvent extraction, or by removing the solvent, partitioning the residue in a mixture of water and a suitable organic solvent and work-up of the organic phase to afford the product.\nThe substituted 2-phenylpyridines of the formula I may contain one or more chiral centers, in which case they are usually obtained as enantiomer or diastereomer mixtures. If desired, these mixtures can be separated into substantially pure isomers using the customary methods for this purpose, such as crystallization or chromatography, including chromatography over an optically active adsorbate. Pure optically active isomers can also be prepared, for example, from suitable optically active starting materials.\nThose substituted 2-phenylpyridines I where R6, R7 and R9=hydrogen can be converted in a manner known per se into their salts, preferably into their alkali metal salts.\nSalts of I where the metal ion is not an alkali metal ion can be prepared by cation exchange of the corresponding alkali metal salt in a conventional manner, similarly ammonium, phosphonium, sulfonium and sulfoxonium salts by means of ammonia, phosphonium, sulfonium or sulfoxonium hydroxides.\nThe compounds I and their agriculturally useful salts are suitable, both in the form of isomer mixtures and in the form of the pure isomers, as herbicides. The herbicidal compositions comprising I control vegetation on non-crop areas very efficiently, especially at high rates of application. They act against broad-leaved weeds and grass weeds in crops such as wheat, rice, maize, soya and cotton without causing any significant damage to the crop plants. This effect is mainly observed at low rates of application.\nDepending on the application method in question, the compounds I, or herbicidal compositions comprising them, can additionally be employed in a further number of crop plants for eliminating undesirable plants. Examples of suitable crops are the following:\nAllium cepa, Ananas comosus, Arachis hypogaea, Asparagus officinalis, Beta vulgaris spec. altissima, Beta vulgaris spec. rapa, Brassica napus var. napus, Brassica napus var. napobrassica, Brassica rapa var. silvestris, Camellia sinensis, Carthamus tinctorius, Carya illinoinensis, Citrus limon, Citrus sinensis, Coffea arabica (Coffea canephora, Coffea liberica), Cucumis sativus, Cynodon dactylon, Daucus carota, Elaeis guineensis, Fragaria vesca, Glycine max, Gossypium hirsutum, (Gossypium arboreum, Gossypium herbaceum, Gossypium vitifolium), Helianthus annuus, Hevea brasiliensis, Hordeum vulgare, Humulus lupulus, Ipomoea batatas, Juglans regia, Lens culinaris, Linum usitatissimum, Lycopersicon lycopersicum, Malus spec., Manihot esculenta, Medicago sativa, Musa spec., Nicotiana tabacum (N. rustica), Olea europaea, Oryza sativa, Phaseolus lunatus, Phaseolus vulgaris, Picea abies, Pinus spec., Pisum sativum, Prunus avium, Prunus persica, Pyrus communis, Ribes sylvestre, Ricinus communis, Saccharum officinarum, Secale cereale, Solanum tuberosum, Sorghum bicolor (s. vulgare), Theobroma cacao, Trifolium pratense, Triticum aestivum, Triticum durum, Vicia faba, Vitis vinifera and Zea mays. \nIn addition, the compounds I may also be used in crops which tolerate the action of herbicides owing to breeding, including genetic engineering methods.\nMoreover, the substituted 2-phenylpyridines I are also suitable for the desiccation and/or defoliation of plants.\nAs desiccants, they are suitable, in particular, for desiccating the aerial parts of crop plants such as potatoes, oilseed rape, sunflowers and soybeans. This allows completely mechanical harvesting of these important crop plants.\nAlso of economic interest is the facilitation of harvesting, which can be achieved by concentrating into a short period of time fruit drop, or reduction of the adherence to the tree, in citrus fruit, olives or other species and varieties of pomaceous fruit, stone fruit and nuts. The same mechanism, i.e. promotion of the formation of abscission tissue between fruit or leaf and shoot of the plant, is also important for readily controllable defoliation of useful plants, in particular cotton.\nMoreover, shortening the period within which the individual cotton plants mature results in improved fiber quality after harvesting.\nThe compounds I, or the compositions comprising them, can be used for example in the form of ready-to-spray aqueous solutions, powders, suspensions, also highly-concentrated aqueous, oily or other suspensions or dispersions, emulsions, oil dispersions, pastes, dusts, materials for broadcasting, or granules, by means of spraying, atomizing, dusting, spreading or watering. The use forms depend on the intended purpose; in any case, they should guarantee the finest possible distribution of the active ingredients according to the invention.\nSuitable as inert auxiliaries are essentially the following: mineral oil fractions of medium to high boiling point, such as kerosene and diesel oil, furthermore coaltar oils and oils of vegetable or animal origin, aliphatic, cyclic and aromatic hydrocarbons, eg. paraffin, tetrahydronaphthalene, alkylated naphthalenes and their derivatives, alkylated benzenes and their derivatives, alcohols such as methanol, ethanol, propanol, butanol and cyclohexanol, ketones such as cyclohexanone, strongly polar solvents, eg. amines such as N-methylpyrrolidone, and water.\nAqueous use forms can be prepared from emulsion concentrates, suspensions, pastes, wettable powders or water-dispersible granules by adding water. To prepare emulsions, pastes or oil dispersions, the 2-phenylpyridines I, either as such or dissolved in an oil or solvent, can be homogenized in water by means of a wetting agent, tackifier, dispersant or emulsifier. Alternatively, it is possible to prepare concentrates comprising active ingredient, wetting agent, tackifier, dispersant or emulsifier and, if desired, solvent or oil, which are suitable for dilution with water.\nSuitable surfactants are the alkali metal salts, alkaline earth metal salts and ammonium salts of aromatic sulfonic acids, eg. ligno-, phenol-, naphthalene- and dibutylnaphthalenesulfonic acid, and of fatty acids, alkyl- and alkylarylsulfonates, alkyl sulfates, lauryl ether sulfates and fatty alcohol sulfates, and salts of sulfated hexa-, hepta- and octadecanols, and also of fatty alcohol glycol ethers, condensates of sulfonated naphthalene and its derivatives with formaldehyde, condensates of naphthalene or of the naphthalenesulfonic acids with phenol and formaldehyde, polyoxyethylene octylphenol ether, ethoxylated isooctyl-, octyl- or nonylphenol, alkylphenyl or tributylphenyl polyglycol ether, alkylaryl polyether alcohols, isotridecyl alcohol, fatty alcohol/ethylene oxide condensates, ethoxylated castor oil, polyoxyethylene alkyl ethers or polyoxypropylene alkyl ethers, lauryl alcohol polyglycol ether acetate, sorbitol esters, lignin-sulfite waste liquors or methylcellulose.\nPowders, materials for broadcasting and dusts can be prepared by mixing or grinding the active ingredients together with a solid carrier.\nGranules, eg. coated granules, impregnated granules and homogeneous granules, can be prepared by binding the active ingredients to solid carriers. Solid carriers are mineral earths such as silicas, silica gels, silicates, talc, kaolin, limestone, lime, chalk, bole, loess, clay, dolomite, diatomaceous earth, calcium sulfate, magnesium sulfate and magnesium oxide, ground synthetic materials, fertilizers such as ammonium sulfate, ammonium phosphate, ammonium nitrate and ureas, and products of vegetable origin, such as cereal meal, tree bark meal, wood meal and nutshell meal, cellulose powders, or other solid carriers.\nThe concentrations of the active ingredients I in the ready-to-use preparations can be varied within wide ranges. In general, the formulations comprise approximately from 0.001 to 98% by weight, preferably 0.01 to 95% by weight of at least one active compound. The active compounds are employed in a purity of from 90% to 100%, preferably 95% to 100% (according to NMR spectrum). The compounds I according to the invention can be formulated for example as follows:\nI 20 parts by weight of the active ingredient of Example No. 9 are dissolved in a mixture composed of 80 parts by weight of alkylated benzene, 10 parts by weight of the adduct of from 8 to 10 mol of ethylene oxide to 1 mol of oleic acid N-monoethanolamide, 5 parts by weight of calcium dodecylbenzenesulfonate and 5 parts by weight of the adduct of 40 mol of ethylene oxide to 1 mol of castor oil. Pouring the solution into 100,000 parts by weight of water and finely distributing it therein gives an aqueous dispersion which comprises 0.02% by weight of the active ingredient.\nII 20 parts by weight of the active ingredient of Example No. 8 are dissolved in a mixture composed of 40 parts by weight of cyclohexanone, 30 parts by weight of isobutanol, 20 parts by weight of the adduct of 7 mol of ethylene oxide to 1 mol of isooctylphenol and 10 parts by weight of the adduct of 40 mol of ethylene oxide to 1 mol of castor oil. Pouring the solution into 100,000 parts by weight of water and finely distributing it therein gives an aqueous dispersion which comprises 0.02% by weight of the active ingredient.\nIII 20 parts by weight of the active ingredient of Example No. 6 are dissolved in a mixture composed of 25 parts by weight of cyclohexanone, 65 parts by weight of a mineral oil fraction of boiling point 210 to 280xc2x0 C. and 10 parts by weight of the adduct of 40 mol of ethylene oxide to 1 mol of castor oil. Pouring the solution into 100,000 parts by weight of water and finely distributing it therein gives an aqueous dispersion which comprises 0.02% by weight of the active ingredient.\nIV 20 parts by weight of the active ingredient of Example No. 9 are mixed thoroughly with 3 parts by weight of sodium diisobutylnaphthalenesulfonate, 17 parts by weight of the sodium salt of a lignosulfonic acid from a sulfite waste liquor and 60 parts by weight of pulverulent silica gel, and the mixture is ground in a hammer mill. Finely distributing the mixture in 20,000 parts by weight of water gives a spray mixture which comprises 0.1% by weight of the active ingredient.\nV 3 parts by weight of the active ingredient of Example No. 14 are mixed with 97 parts by weight of finely divided kaolin. This gives a dust which comprises 3% by weight of active ingredient.\nVI 20 parts by weight of the active ingredient of Example No. 8 are mixed intimately with 2 parts by weight of calcium dodecylbenzenesulfonate, 8 parts by weight of fatty alcohol polyglycol ether, 2 parts by weight of the sodium salt of a phenol/urea/formaldehyde condensate and 68 parts by weight of a paraffinic mineral oil. This gives a stable oily dispersion.\nVII 1 part by weight of the active ingredient of Example No. 6 is dissolved in a mixture composed of 70 parts by weight of cyclohexanone, 20 parts by weight of ethoxylated isooctylphenol and 10 parts by weight of ethoxylated castor oil. This gives a stable emulsion concentrate.\nVIII 1 part by weight of the active ingredient No. 14 is dissolved in a mixture composed of 80 parts by weight of cyclohexanone and 20 parts by weight of Wettol(copyright) EM 31 (=nonionic emulsifier based on ethoxylated castor oil). This gives a stable emulsion concentrate.\nThe herbicidal compositions or the active ingredients can be applied pre- or post-emergence. If the active ingredients are less well tolerated by certain crop plants, application techniques may be used in which the herbicidal compositions are sprayed, with the aid of the spraying equipment, in such a way that as far as possible they do not come into contact with the leaves of the sensitive crop plants, while the active ingredients reach the leaves of undesirable plants growing underneath, or the bare soil surface (post-directed, lay-by).\nThe rates of application of active ingredient are from 0.0005 to 3.0, preferably 0.0005 to 1.0, kg/ha of active substance (a.s.), depending on the control target, the season, the target plants and the growth stage.\nTo widen the spectrum of action and to achieve synergistic effects, the 2-phenylpyridines I may be mixed with a large number of representatives of other herbicidal or growth-regulating active ingredients and then applied concomitantly. Suitable components for mixtures are, for example, 1,2,4-thiadiazoles, 1,3,4-thiadiazoles, amides, aminophosphoric acid and its derivatives, aminotriazoles, anilides, (het)aryloxyalkanoic acids and their derivatives, benzoic acid and its derivatives, benzothiadiazinones, 2-aroyl-1,3-cyclohexanediones, hetaryl aryl ketones, benzylisoxazolidinones, meta-CF3-phenyl derivatives, carbamates, quinolinic acid and its derivatives, chloroacetanilides, cyclohexane-1,3-dione derivatives, diazines, dichloropropionic acid and its derivatives, dihydrobenzofurans, dihydrofuran-3-ones, dinitroanilines, dinitrophenols, diphenyl ethers, dipyridyls, halocarboxylic acids and their derivatives, ureas, 3-phenyluracils, imidazoles, imidazolinones, N-phenyl-3,4,5,6-tetrahydrophthalimides, oxadiazoles, oxiranes, phenols, aryloxy- and hetaryloxyphenoxypropionic esters, phenylacetic acid and its derivatives, phenylpropionic acid and its derivatives, pyrazoles, phenylpyrazoles, pyridazines, pyridinecarboxylic acid and its derivatives, pyrimidyl ethers, sulfonamides, sulfonylureas, triazines, triazinones, triazolinones, triazolecarboxamides and uracils.\nIt may furthermore be advantageous to apply the compounds I, alone or in combination with other herbicides, in the form of a mixture with other crop protection agents, for example together with agents for controlling pests or phytopathogenic fungi or bacteria. Also of interest is the miscibility with mineral salt solutions, which are employed for treating nutritional and trace element deficiencies. Non-phytotoxic oils and oil concentrates may also be added."}
{"text": "The present invention relates generally to a tachometer used to detect engine speed, and more particularly, to a universal tachometer operative to measure the speeds of engines with different numbers of cylinders.\nTachometers have been broadly applied to measure the angular velocity, typically measured in units of revolutions per unit (RPM). Among various tachometers, the contactless magnetic pickup tachometers are particularly popular for measuring the engine speed of automobiles. For measuring the engine speed, the magnetic pickup tachometer is typically disposed near an ignition timing wheel made of ferromagnetic material. When each tooth of the ignition timing wheel passes through the magnetic pickup tachometer, a pulse is generated. Normally, either the period of the pulse is measured, or the number of pulses per unit time is counted. According to the number of the teeth for the ignition timing wheel and the measured period or counted pulse, the speed can be obtained.\nAs the above tachometer has to be installed near the crankshaft of the engine, the installment is difficult and laborious. Therefore, companies have developed tachometers with inductor pickup coupled to coil wires of the engine, avoiding the need to measure the magnetic field induced by the rotation of the crankshaft. However, as the coil wires are not readily accessible in the engines in many new vehicles, the installation may remain difficult.\nFurther, conventional tachometers are designed for the engine with a specific number of cylinders. To measure the speed of engines with different numbers of cylinders, different tachometers are required. Moreover, in many newer engines, two sparks are generated for each operation cycle for each cylinder. Using the conventional tachometer may incorrectly read the second spark as a separate cycle, doubling the reported engine speed."}
{"text": "The present invention relates to an X-ray beam device for X-ray analytical applications such as X-ray Diffractometry, High Resolution X-Ray Diffractometry, X-ray Reflectometry comprising an X-ray Source, typically a point focus X-ray source, and an x-ray focusing optical assembly for focusing monochromatic x-rays on a sample.\nThe invention will be more particularly adapted for High Resolution X-ray Diffraction applications. High Resolution X-ray Diffraction is a well established technique for carrying out the analysis of high quality thin layer of near perfect single crystals materials.\nHigh resolution diffraction analysis requires to illuminate the sample with an X-ray beam having a very small divergence and a high energetic resolution. This is necessary to determine the parameters such as lattice spacing with a sufficient precision. Typically for High Resolution X-ray Diffraction applications in the case where Copper radiation is used, the energy resolution required is such that the incident x-ray beam needs to be filtered so that only Kα1 line is transmitted to the sample. The Kα2 line from the Copper Kα doublet has consequently to be filtered.\nThe X-ray Beam Device according to the invention typically comprises a high brilliance X-ray source which can be either a sealed tube X-ray source or a rotating anode X-ray source with a point focus geometry. Laboratory X-ray sources have isotropic radiations and the characteristic X-ray spectrum emitted by an X-ray source with a Copper target is composed of several parasitic X-ray peaks and bremsstrahlung x-rays. Moreover the Kα1 and Kα2 lines are merged in the Kα doublet.\nX-ray multilayer optical elements composed of a Bragg reflective structure can be used to collect the beam emitted by an x-ray source in order to monochromatize and adapt the x-ray beam in real space and in angular space toward a sample. However the use of a multilayer optical element to filter the X-ray beam emitted by a source will be unsufficient for high resolution diffraction applications as for Copper X-ray beam radiation the Kα1 and Kα2 lines are merged after the reflection on the multilayer optical element.\nNatural crystals (e.g. Si, Ge, LiF) can be used as X-ray optical elements in order to monochromatize an x-ray beam to the level of energetic resolution required for High Resolution X-ray Diffraction. However the efficiency of such natural crystals is limited when curved crystal elements necessary for focusing the X-ray beam emitted from the source are used (curvature deteriorates the very small lattice spacing structure). Such natural crystals typically have a corresponding angular acceptance of 10 to 30 Arcseconds. FIG. 1 is a schematic graph illustrating an example of angular acceptance ΔθM of a crystal monochromator in terms of reflectivity.\nOne optical system arrangement typically used for high resolution x-ray diffraction applications is composed of an x-ray multilayer conditioning optic that is used to convert the divergent X-ray beam emitted by the X-ray source into a parallel beam before entering a crystal monochromator arranged downstream the multilayer optic. Such multilayer conditioning optic will be referred to a collimating optic in the description of the invention. The collimating effect of a multilayer conditioning optic is referring to the optical effect where a divergent beam collected by the optic is converted in a substantially parallel beam having a divergence lower than 1 milliradians typically.\nFIG. 2 is illustrating a type of multilayer conditioning optic 20 producing a collimating effect wherein the optical element comprises a reflective surface being shaped according to two curvatures Cx and Cy corresponding to two different directions which are the sagittal direction (direction Y in FIG. 2) and the meridional direction (direction X in FIG. 2). These directions that will be referred to in the description of the invention can be defined with respect to the general direction of propagation of the X-ray beam:—the meridional direction being the mean direction of propagation of the x-ray beam (and more precisely the mean direction between the mean direction of propagation of the beam before and after its reflection on the optical assemblies concerned), and the sagittal direction being the horizontal transverse direction of this meridional direction (the vertical being defined by the mean normal to the part of the reflective surface of the optical assemblies which will be described and used for reflecting the incident x-ray beam). The incoming x-ray beam is after a single reflection on the optic 20 collimated in two-dimensions which are the sagittal plane (a sagittal plane being the plane defined by the sagittal direction and by the mean direction of propagation of the x-ray beam at the exit or at the entry of the concerned optical element) and in the meridional plane (the meridional plane being the plane defined by the meridional direction and the mean normal to the part of the reflective surface of the optical assemblies which will be described and used for reflecting the incident x-ray beam).\nSuch optical system arrangement with a collimating multilayer optic 20 arranged upstream a crystal monochromator is efficient for applications requiring a few hundred microns or a millimiter sized x-ray beam at the sample position. The X-ray spot size at the sample position as it is defined in the invention is given by the dimensions Ws and WT of the beam in the two directions perpendicular to the general direction of propagation of the beam. As an indication, Ws and WT dimensions are illustrated in FIG. 2 for the x-ray beam at the exit of the two-dimensional collimating mirror 20.\nHowever such collimating optical systems are limited for applications requiring a small spot dimension at the sample position (in the order of few hundred microns or less than 100 microns). Indeed the flux will be reduced proportionally to the surface of the spot.\nAnother known optical system arrangement is to use a two-dimensional focusing optic with a multilayer coating adapted upstream the crystal monochromator in order to collect the divergent x-ray beam emitted from the source and to focus such beam on a small spot at the sample. However, some intensity is lost due to the divergence of the X-ray beam incoming the monochromator in the scattering plane of the monochromator (the scattering plane of the monochromator is the plane including the incident beam and the diffracted beam) resulting from the focusing effect of the upstream multilayer optic.\nFIG. 3 illustrates an optical arrangement known from state of the art where an optical element 21 with a multilayer coating collects an incident X-ray beam X1 emitted by a point source 10 and conditions such beam towards a crystal monochromator M by collimating the beam X2 in one dimension and focusing the beam X2 in the other dimension. As illustrated in FIG. 3, the x-ray beam X2 outcoming from the optical element 21 is collimated in the meridional plane of the optical element 21 corresponding to the scattering plane of the monochromator, and is focused in the sagittal plane of the optical element 21 corresponding to the sagittal plane of the Monochromator. The focusing in the sagittal plane of the Monochromator M enables to concentrate a higher intensity beam on a reduced spot size at the sample 40. This focusing effect, leading to a slightly divergent beam in the sagittal plane is possible due to a larger tolerance of the crystal monochromator in this plane (compared to the tolerance on beam divergence in the scattering plane which is the angular acceptance as shown in FIG. 1). This difference of tolerances on incident beam divergence for a crystal monochromator is referenced for example in the Document entitled “Parallel Beam Coupling into channel-cut monochromators using curved graded multilayers” from M. Schuster and H. Gobel published in J. Phys. D: Applied Physics 28 (1995) A270-275 (a maximum divergence value of 1.95° in the sagittal plane is given as an example for a Germanium(022) crystal monochromator for the Copper Kα1 line).\nHowever, the overall usefull flux of such optical arrangement that is illustrated in FIG. 3 is limited due to a reduced divergence of the x-ray beam diffracted by the crystal in the scattering plane of the crystal monochromator. Indeed the divergence of the X-ray beam in this scattering plane of the monochromator will be limited to the reflection width of the crystal ΔθM which is typically between 10 and 30 Arcseconds. As it will be illustrated later in the description of the invention, even for High Resolution X-ray Diffraction applications it can be accepted to have a higher divergence of the incidence X-ray beam that is illuminating the sample compared to the angular acceptance ΔθM of a traditional crystal monochromator. However it is not possible to increase the angular acceptance of a traditional crystal monochromator while keeping a high energetic resolution.\nAn object of the present invention is to increase the usefull flux compared to traditional x-ray beam device for High Resolution X-ray Diffraction Applications in particular for applications requiring a small spot (in the order of few hundred microns or less than 100 microns).\nA further object of the present invention is to achieve very small spot dimensions on the sample, in particular spot dimensions smaller than 50 microns for X-ray analytical applications requiring such high spatial resolution as in semiconductor metrology applications."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an imaging screen for electrophoresis applications and, more particularly, to an imaging screen for detecting and recording the impingement of beta particles or visible light emitted from electrophoresis gels.\n2. Description of the Related Art\nIn order to sequence deoxyribonucleic acid (DNA), it is necessary to chemically treat and label DNA fragments prior to performing electrophoresis. Specifically terminated fragments of the DNA are placed in wells in a gel and an electric field applied, resulting in the motion along well-defined lanes in the gel. Depending on the length of the fragments and the time for performing the electrophoresis, they will travel more or less far along the lanes. The result of the electrophoresis operation is a pattern of bands of DNA fragments.\nIf the fragments have been labelled with a beta particle emitter, it is possible to obtain a two-dimensional map of the pattern by placing X-ray film in relatively intimate contact to the gel. The X-ray film is somewhat sensitive to the beta particle radiation, and generally after long exposure will record the pattern. Unfortunately, such a procedure typically requires a time period on the order of days for the film to register a sufficient number of emitted beta particles to provide an accurate representation of the pattern of bands.\nAnother drawback to this method is associated with the properties of X-ray film. Because of its limited dynamic range, and the complexity of obtaining quantitative information using such film, efforts have been undertaken to develop simpler, speedier and more accurate methodologies to detect radioactively labelled DNA fragments in electrophoresis gels.\nFour beta emitters are typically employed for tagging DNA fragments--i.e., .sup.32 P, .sup.35 S, .sup.14 C, and .sup.3 H. The beta particles end point energies are 1.71 MeV, 168 keV, 156 keV, and 18.6 keV, respectively. The electron ranges vary greatly, with .sup.32 P having the longest range, and .sup.3 H having the shortest range.\nRecently developed devices for obtaining quantitative information from radioactively labelled DNA, RNA or protein fragments in electrophoresis gels generally utilize two coincident grids which sense the beta emissions from the gels. Such a device is marketed as the Betascope 603 Blot Analyzer. That device can perform an analysis of the emitted beta particles in a number of hours. However, the Betascope device, and similar such devices, cannot readily detect the shorter range and lower energy beta particles. In addition, its spatial resolution is poor.\nFuji has developed a BaFBr imaging screen system which can work with .sup.32 P, .sub.35 S, and .sup.14 C, but is at least 1-2 orders of magnitude slower then the Betascope Blot Analyzer in processing the information contained in the radioactively labelled DNA, RNA or protein carrying electrophoresis gel.\nAccordingly, there is a need for a simple, speedy and accurate methodology for detecting radioactive tracers emitted from electrophoresis gels. It would be additionally advantageous if such methodology could also be utilized to detect dye-tagged and chemiluminescent tracers which emit visible light, which the Fuji BaFBr system cannot detect."}
{"text": "This invention is generally referred to a methodology of structural prognostic, and more specifically to a method and a system for the prognostic of a mechanical structure, in particular a structure of aircraft, adapted to evaluate the presence or to monitor the rise-up of damages or defects induced in a structure by operative loads and/or service events.\nIn the maintenance methodologies of systems (parts of products or complex products) it is extremely interesting the possibility of reducing unexpected failures, through the monitoring of some indicators of the system status.\nAccording with the known technique, in aeronautics the rise-up of damages or defects in a mechanical structure of an aircraft, such as metallic or composite structure, e.g. the fuselage or the wing structure, are prognosticated through a historical recording of the events among which events that caused damages due to accidental impacts during production (impact of a tool) or service (impact with hail or rubble) and the loads born by the structure, or by an estimate of the fatigue underwent by the structure, starting from the knowledge of its characteristics of mechanical strength under typical loads in service conditions.\nIn particular, for composite structures, accidental impacts produce effects scarcely visible from outside, but they may cause relevant damages within the structure (e.g. delamination).\nThis technique is however expensive and not very accurate, because it does not indicate in real time the changes and the physical and mechanical conditions of the monitored structure."}
{"text": "The present disclosure relates to a method, non-transitory computer readable recording medium, and apparatus for scheduling virtual machine monitor, and more particularly, relates to a method, non-transitory computer readable recording medium, and apparatus for scheduling virtual machine monitor for reducing an input/output (I/O) event delay of scheduling.\nVirtualization technology refers to the technology of efficiently using a resource by integrating a plurality of physical resources into a logical resource, or on the contrary, dividing a single physical resource into a plurality of logical resources. The virtualization technology may operate a plurality of systems in a single machine by introducing a virtualization layer. The software of abstracting the physical resource is referred to as a virtual machine monitor or a hypervisor, and there are well-known virtual machine monitors such as Hyper-V, VirtualBox, VMware, Xen, Kernel-based Virtual Machine (KVM), etc.\nThe virtual machine refers a logical resource abstracted by a virtual machine monitor, that is, the virtual machine does not mean an actual machine but means a machine environment in which a physical machine is likely to exist, and in Xen, the virtual machine may be referred to as a domain. A relationship between the virtual machine monitor and the virtual machine is similar to a relationship between an operating system and a processor. The operating system provides an abstracted physical resource to a plurality of processors by abstracting the physical resource, and similarly, the virtual machine monitor provides virtualized physical resources to the virtual machine by virtualizing a central processing unit (CPU), a memory, and an input/output (I/O).\nThe virtualized CPU (VCPU) refers to the virtual machine monitor logically abstracting the CPU which is an actual physical resource, and a plurality of virtual machines operating in a single physical machine have not the actual CPU but the VCPU. The VCPU corresponds to the processor in the relationship between the operating system and processor, and the virtual machine monitor also selects the VCPU which is to use the CPU resource through a scheduling policy as if the operating system selects the processor which is to use the CPU resource. As such, scheduling of determining the VCPU which is to use the CPU resource is referred to as virtual machine scheduling.\nA credit scheduler is one of the virtual machine scheduling methods of Xen which is a Linux open resource, and uses a boost mechanism in order to improve event responsiveness. The boost mechanism changes a state of a corresponding VCPU into an executable state when an event is delayed in the VCPU which is in a block state, and allocates BOOST which is the highest priority. The corresponding VCPU is in a block state when there is no executable task in a domain or the VCPU holds a specific value or more credit since a credit is not consumed, and is not executed until the event is generated. The VCPU having BOOST which is the highest priority is preferentially executed, and improves the responsiveness of the domain by rapidly processing the delayed event.\nHowever, since the VCPU performing a CPU-intensive operation continuously consumes the credit and a probability in which an executable task exists in the domain is high, the VCPU cannot receive a benefit of a conventional boost mechanism. Further, when there are a plurality of virtual machines which are in an active state, a virtual machine scheduling delay time is increased, and event responsiveness is further decreased. In a virtualization environment, since an interrupt generating in the I/O operation of the virtual machine is processed using an event mechanism, I/O performance as well as the I/O responsiveness may be decreased when the event responsiveness is decreased."}
{"text": "In optical communication systems, the need may arise to multiplex different wavelength signals onto a single fiber. For example, wavelength multiplexing is one method for achieving full bidirectional transmission on a single fiber. In its simplest form, a bidirectional system may comprise two stations, S1 and S2, which transmit information at wavelengths λ1 and λ2, respectively. Thus, station S1 needs a transmitter, which operates at λ1 and a receiver, which is tuned to wavelength λ2. Station S2, obviously, has the opposite requirements. Each station also needs a duplexing element to inject both wavelengths onto the single transmitting fiber. Although simple in theory, such an arrangement is cumbersome in implementation. For example, each station comprises a separate transmitter, receiver, and duplexer. Therefore, some sort of optical coupling must also be provided, for example, by using optical waveguides. Such coupling requires many expensive and time-consuming adjustments to achieve optimum alignment. Additionally, the optical losses attributed to this coupling, including attachment between the duplexer, fiber, transmitter and receiver, may degrade the overall performance of the station to an unacceptable level.\nAn alternative to this straightforward implementation is disclosed in U.S. Pat. No. 4,592,619 issued to E. Weidel on Jun. 3, 1986. Weidel discloses an optical coupling element utilizing a variety of microoptic elements with spherical and plane surfaces for collimating, focusing and redirecting transmitted/received light waves. Although an improvement over the prior art, the Weidel arrangement utilizes at least one optical element, which must be traversed twice by a received light signal. Further, Weidel is necessarily limited to providing coupling between both a transmitter and receiver to an optical fiber. However, there exist situations wherein a pair of transmitters, operating at different wavelengths, must be coupled over the same fiber (unidirectional transmitter).\nThus, a need remains in the prior art for a dual wavelength optical coupler, which is robust in design and is capable of operating in either a bidirectional mode (transmitter and receiver) or unidirectional mode (two transmitters or two receivers).\nAn attempt to solve the problems that occur in the system of U.S. Pat. No. 4,592,619 is made in U.S. Pat. No. 4,904,043 issued in 1990 to R. Schweizer. This technique was developed in AT&T Bell Laboratories. In this patent, Schweizer describes a device, which practically is the nearest prototype of modern bidirectional transceivers as it contains all elements that found use in subsequent devices of this type. Dual wavelength coupling is achieved utilizing a set of three lenses and a dichroic filter, all held in a precision die-cast housing with the active devices. In one embodiment, the coupler may be used as a bidirectional transceiving device that includes a LED operating on a first wavelength and a PIN receptive to a different wavelength. In another embodiment, the coupler may be used as a unidirectional device, including either a pair of LEDs at different wavelengths or a pair of PINs at different wavelengths. One of the objects of the design described by Schweizer is to avoid active alignment of the components forming the coupler. By careful choice of the lenses, alignment tolerances may be minimized to the extent that the filter and lenses may be merely placed in their proper locations within the housing. Another aspect is to provide a coupler design, which is flexible enough to be utilized with a number of different lenses, as well as different transmitting and receiving wavelengths.\nHowever, a main disadvantage of the device disclosed in U.S. Pat. No. 4,904,043 consists in that a photoreceiver should always be optically coaxial with an optical axis of one of two light sources that generate light of a working wavelength λ1 or λ2. Such a design is inconvenient for suppression of crosstalk, e.g., a parasitic signal with the wavelength λ2 when a photodiode receives a useful signal with wavelength λ1.\nAll further developments in the field of bidirectional transceivers had design close to the aforementioned device developed by AT&T Bell Laboratories, but with positions of the photodiode and the second light source reversed for obviating the inconvenience inherent in the device described in U.S. Pat. No. 4,904,043.\nU.S. Pat. No. 5,485,538 issued on Jan. 16, 1996 to T. Bowen et al. discloses a typical bidirectional transceiver with orthogonal arrangement of two light beams with wavelength λ1 and λ2. The device comprises a compact optical transceiver that includes a ceramic mounting block with a laser diode abutting a first end of the mounting block for generating light of a first wavelength. A holographic optical element (HOE) is positioned adjacent a laser diode and acts as a hologram lens which receives and focuses the generated light to the end face of an optical fiber which is attached to the compact optical transceiver. A glass element is mounted on the mounting block between the diode and the optical fiber end face and includes a dichroic beam splitter that passes light of a first wavelength λ1 and deflects light of a second wavelength λ2. The dichroic beam splitter is mounted in an angular positioning groove of the mounting block and receives and passes the generated light of the first wavelength λ1 which has been focused by the HOE. From a remote transmitter, light of a second wavelength λ2 is then transmitted through the optical fiber to the compact optical transceiver from a direction opposite to that of the light generated by the laser diode. The light transmitted from the optical fiber is then transmitted and passed through a section of the fiber supported by a ferrule attachable to the optical fiber and attached in a second V-shaped positioning groove of the ceramic mounting block. The light of the second wavelength is sent to the beam-splitter and is deflected through a bore hole in the ceramic mounting block to a detector abutting the bottom of the mounting block. A blocking filter can be included for blocking light of wavelengths other than the first and second wavelengths from the detector.\nOne disadvantage of the transceiver of U.S. Pat. No. 5,485,538 consists in that all optical elements, including HOE, used for beam management transform the beams into converging or diverging beams. This creates significant problems for alignment of the beams because all the optical elements must be adjusted simultaneously. Another disadvantage, which is inherent in all orthogonal bidirectional transceivers, consists in that their geometry is unsuitable for use in conjunction with high-speed controllers.\nAnother U.S. Pat. No. 5,487,124 issued on Jan. 23, 1996 to the same applicants as the previous patent describes a bidirectional transceiver that differs from the one described in U.S. Pat. No. 5,485,538 by the fact that the HOE was replaced by a GREEN lens. It can be clearly seen from FIG. 2 of U.S. Pat. No. 5,487,124 that the beams of the light source are transformed into converging and diverging beams. Therefore the device of this patent entails all disadvantages of the previously described design, including speed limitations due to geometry.\nIt should be noted that both designes disclosed in two previous patents were developed at the Whitaker Corporation, Wilmington, Del. In an attempt to solve problems associated with difficulties of alignment mentioned in two previous patents, the Whitaker Corporation developed a new design described in U.S. Pat. No. 5,621,573 issued on Apr. 15, 1997 to W. Lewis, et al. This patent describes a bidirectional link that allows sequential or simultaneous transmission and reception of optical signals using conventional components. To effect the relatively simple alignment of the devices and components, the emitter or transmitter is disposed in a subassembly having a sub-housing with the required optical focusing beam splitting elements disposed therein. This subassembly is optically aligned in a relatively simple active alignment process, and the subassembly is then inserted into the main housing of the bidirectional link. At this stage, a detector is mounted in the main housing and is aligned optically with the pre-aligned elements of the sub-assembly described above. The detector is then fixed using common adhesive and the assembly of the bidirectional link is complete. In this device, the inventors for the first time used so-called canted fibers and lens subassemblies that comprise a lens and a fiber preassembled in a common ferrule-type sub-housing that simplifies manipulation and adjustment. The fiber end face can be slightly inclined from a perpendicular to the optical axis of the lens for decreasing reflection and thus for improving optical coupling. Although the bidirectional transceiver of U.S. Pat. No. 5,621,573 partially simplifies the alignment procedure due to the use of preassembled units, the limitations by speed remain unsolved due to the use of practically the same geometry as in all previously known orthogonal arrangements.\nIn subsequent years, designs of orthogonal bidirectional transceivers were improved with a new technique developed by Lucent Technologies Inc. This new technique was aimed at improved alignment, more efficient optical coupling, and suppression of crosstalk. Thus, U.S. Pat. No. 5,796,899 issued to T. Butrie et al. on Aug. 18, 1998 describes an optical transceiver assembly for use in a bidirectional system that includes a beam splitter to direct an incoming signal to a photodiode. An outgoing signal from a laser diode is partially transmitted and partially reflected by the splitter. The reflected signal, which may reach the photodiode, constitutes crosstalk which is reduced by means of a cavity positioned to receive the reflected signal and an oblique surface within the cavity adapted to prevent much of the reflected signal from reaching the photodiode.\nU.S. Pat. No. 5,838,859 issued on Nov. 17, 1998 to the same applicants as in U.S. Pat. No. 5,796,899 describes an optical transceiver assembly for use in a bidirectional system that includes a beam splitter to direct an incoming signal to a photodiode. An outgoing signal from a laser diode is partially transmitted and partially reflected by the splitter. The reflected signal, which may reach the photodiode, constitutes crosstalk, which is reduced by orienting the polarization direction of the splitter essentially parallel to that of the outgoing signal from the laser diode. In another embodiment, which enhances coupling efficiency, a single element aspheric lens is positioned between the laser diode and the splitter.\nAnother device developed by Lucent Technologies Inc. is a modular form that improves accuracy of alignment and makes the device suitable for mass production. This device is described in U.S. Pat. No. 5,841,562 issued on Nov. 24, 1998 to S. Rangwala, et al. In accordance with one aspect of this invention, a transceiver comprises a transmitter module and a receiver-splitter module, which is plugged into a self-aligning socket of the transmitter module. In one embodiment, the transmitter module includes a light source lensed to an opening in the socket, and the receiver-splitter module includes a ferrule, which is plugged into the socket. The ferrule carries an optical fiber so that one end of the fiber is optically coupled to the light source. This coupling enables an outgoing optical signal to be partially transmitted to a fiber pigtail located at the opposite end of the ferrule. A splitter is located at the other end of the fiber so that an incoming optical signal on the fiber pigtail is partially reflected to a light detector.\nU.S. Pat. No. 6,075,635 issued on Jun. 13, 2000 to T. Butrie at al. describes a bidirectional optical transceiver developed at Lucent Technologies Inc. that includes a housing in which a light source, lens, beam splitter, photodetector and an optical fiber are mounted. The lens focuses an outgoing optical signal from the source through the splitter to the end face of the fiber. The splitter directs an incoming optical signal to the photodetector. In order to reduce reflections of the outgoing signal from the end face of the fiber, and hence crosstalk, without also sacrificing significantly the coupling efficiency to the fiber, the fiber end face is beveled at an acute angle φ to the normal to the common axis of the source, splitter and fiber, and the fiber is tilted at an acute angle θ to the same axis. In a preferred embodiment, which further enhances coupling efficiency, the fiber end face is beveled at an even smaller acute angle φ′ to the normal to the fiber axis, and the fiber axis is tilted at an acute angle θ to the common axis, where φ′ is about 2θ.\nA disadvantage common to all four last-mentioned patents of Lucent Technologies Inc., as well to all preceding structures, is that the orthogonal geometry used in the bidirectional optical transceivers described in the aforementioned patents make it difficult to use such devices in high-speed systems with frequencies of 500 MHz or higher.\nIn fact, all known bidirectional transceivers, as well as those described above, are based on the use of standard commercially-produced laser diodes and photodiodes, such as, e.g., TO CAN packages. The orthogonal arrangement and geometrical dimensions of the aforementioned laser diodes and photodiodes required relatively long lead wires (not less than 10 mm) for commutation with the PC board circuitry. This limitation restricts the speed of transmission data through the bidirectional transceiver.\nIt is known that inductance depends on the length of the conductor. Therefore, if one takes as a LD series resistance 8 Ohms and accepts 1 nH/mm inductance, it is easy to evaluate that the integration time constant of LD (T=L/r) for 10 mm lead wire is greater than 1.25 ns and for 3 mm lead wire is greater than 0.375 ns, correspondingly. The first lead wire limits gigabit applications of the bidirectional transceiver, while the second one allows to operate at data rates coming to 2.5 Gb/s."}
{"text": "(1) Field of Invention\nThe present invention related generally to improvements in the design of forced circulation kilns or chambers for drying and, in particular, to kilns for conditioning or seasoning lumber or similar material.\n(2) Description of the Prior Art\nNewly cut lumber contains varying amounts of moisture which may range from 30% to more than 200% by weight depending on the species of the wood and other factors. The dimensions of the cut wood, as well as its strength, will vary as a function of the moisture content. Thus, in order to ensure greater dimensional stability and strength, as well as lower shipping costs, it is necessary to reduce the moisture content of the cut wood.\nRemoval of the moisture from the cut wood is usually referred to as curing. The two commonly used methods of curing are known as the kiln-dry method and the air-dry plus kiln-dry method. The primary difference between the two methods is the air-dry plus kiln-dry method includes the additional steps of first arranging the lumber on elevated racks in open yards to take advantage of the moisture removing capacity of the natural atmosphere. Depending on the type of wood, method of stacking, and other environmental conditions, between 40 and 300 days may be necessary to reduce the moisture content to 20%.\nThe kiln-dry method may be used with or without the air-dry pre-drying step. The lumber is tiered upon racks or kiln trucks and placed in the kiln where the dry bulb temperature and the wet bulb temperature of the surrounding air are maintained at various predetermined conditions and according to published schedules to reduce the moisture content to an amount near that desired in final use. Typically, the desired moisture content is between 6 and 10%.\nWhile kiln-drying is significantly more rapid and more controllable than simple air-drying, it adds significantly to the cost of the cured lumber. A conventional lumber drying kiln is both costly to construct and to operate. First, the conventional lumber-drying kiln must be large enough to accomodate large amounts of stacked lumber. Second, the walls of the kiln must be insulated in order to prevent excessive heat loss. Finally, electricity for the large capacity fans, steam for heating the air in the kiln, and high maintenance costs due to the moisture-laden atmosphere of the kiln result in high operating costs.\nOne example of such a well-known conventional kiln is shown in U.S. Pat. No. 3,131,034 to Marsh. Such kilns are usually provided with two or more longitudinal rows of vents which serve to alternately admit air to the kiln and to exhaust the kiln atmosphere depending on the direction of the fan rotation within the kiln. Each of these vents usually includes a body member aligned to a complementary metal opening in the top of the kiln. The upper portion of the body member is hingedly connected to a vent cap having operating arms which are actuated in response to the direction of the fan operation. Thus, in operation, outside air is brought in on the inlet side of the fan assembly and blown through the heating elements, then a portion of the heated air is exhausted out of the exhaust vent while the larger portion of the air passes around and through the stacked lumber.\nCertain disadvantages become apparent with such a design. Firstly, a significant amount of treated air passes through the steam heaters and is then exhausted through the outlet vents. Secondly, a significant portion of the total fan airflow goes out through the vents before passing through the stacked lumber. Finally, the large number of vents located in the roof adds considerably to the cost of construction.\nBecause of the high temperature, high humidity in the interior of the drying kiln, it has also been the practice in the past to locate the drive motor outside of the kiln proper (see Marsh esp. FIG. 1). Because the fans and the associated motors are conventionally located above the lumber stacks, it has also become the practice to construct a separate room, isolated from the drying chamber, which may house the motor and associated controls. One example of such a separate control room is shown in U.S. Pat. No. 4,098,008, issued to Schuette et al (see esp. FIG. 4).\nCertain disadvantages become apparent with such a design. First, a completely separate structure must be constructed adjacent to the drying kiln which adds significantly to the cost of the drying kiln. Second, the drive shaft between the fans and motor must be significantly lengthened to reach the secondary building. This adds to the construction cost, maintenance, and probability of damage due to the higher vibration of the longer drive shaft. As a result, this design has not been well received by the industry.\nFinally, in the past during the construction of kilns, each kiln has been custom designed and fabricated. As such, once the construction of the housing of the kiln has progressed to a certain stage, it has been the conventional practice for the various fans, heat exchangers, humidifiers, baffles, etc. to be individually stalled prior to completing construction of the kiln. This results in a tedious and prolonged construction of the kiln such that the normal time for construction is approximately eight weeks. This has been due, in a large part, to the complicated roof designs of the previous kilns.\nIt has thus become desirable to develop a lumber drying kiln which will more effectively utilize the energy necessary for drying lumber. It has also become desirable to develop a lumber kiln structure which will better facilitate access to the fan motors without unnecessary costs and maintenance. Finally, it has become desirable to develop a lumber kiln structure which is particularly adaptable to prefabrication, thereby eliminating the long lead time necessary to construct present day lumber kilns."}
{"text": "Information models based on extensible markup language (XML) may be complex and difficult to implement. This may be particularly evident in industry standard information models (e.g., healthcare) that possess one or more of the following characteristics:                1. Large number of complex types (elements and attributes)        2. Generic/abstract (high-level domain concepts)        3. High degree of extensibility/customization is possible        4. Recursive—leads to instances with arbitrarily deep nesting        5. Instances contain a high ratio of fixed to variable data        \nTo build applications that produce, consume and validate such XML-based information models, additional constraints are often placed on the base standard to define use-case specific structure and content. Software code created to implement these capabilities and validate these constraints may be very specific to a particular application, and can be difficult to maintain and reuse across multiple applications. It is possible to create abstractions on top of these data models and write transformations between the abstractions and the complex schema, but the process is typically manual and labor intensive."}
{"text": "During semiconductor fabrication of memory devices, to save space, the capacitor of a storage cell must reduce in size and yet maintain adequate capacitance to retain a sufficient charge, such as during Dynamic Random Access Memory (DRAM) operation. There are several approaches to the capacitor design, for example trench capacitors formed in the substrate of a wafer or a stacked capacitor formed above the wafer substrate, to name two.\nRegardless of the design chosen, the size of the capacitor must be reduced and yet maintain sufficient capacitance as mentioned previously. Two of the main contributors to capacitance are the surface area of the capacitor plates and the dielectric quality of the insulator separating the capacitor plates. Major engineering efforts have gone into improvements in both areas.\nOnce the capacitor design is chosen, another factor should be considered and that is the physical connection between the capacitor plate and the underlying conductor (i.e., a conductive plug between the capacitor plate and a source/drain of an access transistor or the source/drain of the access transistor itself). It is desirable that the physical connection between the capacitor plate and the underlying conductor consist of as low resistance as possible and thus provide for low “time at temperature” (DT) flow during fabrication processing.\nA significant focus of the present invention is the development of a capacitor having a high surface area within a confined area that possesses a low DT flow. Thus, the present invention comprises a capacitor structure having a high surface area within a confined area that possesses a low DT flow and a method to fabricate same, which will become apparent to those skilled in the art from the following disclosure."}
{"text": "Understanding the reservoir fluids, hydrocarbon-generating potential, and formation characteristics of a reservoir rock formation is important to estimate the reserve and to improve the efficiency of oil field operations such as drilling, well completion, and production. To this end, various types of formation and reservoir analysis are performed using downhole tools. Further, various types of rock sample analysis are performed using laboratory tools at earth's surface. One example of formation or rock sample analysis involves nuclear magnetic resonance (NMR) phenomena.\nAn example NMR tool includes a magnet assembly that produces a static magnetic field (B0), and a coil assembly that generates radio frequency (RF) electromagnetic signals (including a perturbing magnetic field (B1)) and that detects magnetic resonance phenomena. The behavior of nuclei in presence of B0 and B1 has been correlated with formation rock and fluid properties such as the amount of hydrogen in pore space fluids and/or the porosity.\nTotal organic carbon (TOC) is an important reservoir quality indicator, particularly for unconventional source rock reservoirs. It is well-known that TOC can be obtained from laboratory mineralogy and fluid analysis using techniques such as the combination of X-ray diffraction (XRD) and X-ray fluorescence (XRF). However, laboratory tests are expensive and time consuming. Performing such tests for a large number of samples corresponding to a particular reservoir or even a single well is impractical. Nuclear logging based mineralogy analysis can also be used to estimate TOC, provided that all elements in the formation can be resolved, which is not always possible. Combining density and NMR logging measurements can also be used to estimate TOC, with an assumption that NMR logs do not include any signals from organic matters, liquid or solid, in the source rocks, an assumption not always valid.\nIt should be understood, however, that the specific embodiments given in the drawings and detailed description below do not limit the disclosure. On the contrary, they provide the foundation for one of ordinary skill to discern the alternative forms, equivalents, and other modifications that are encompassed in the scope of the appended claims."}
{"text": "To centre a workpiece on a lathe a hole is first drilled at one end of the workpiece and secondly a centre of the lathe is located in the hole so as to rotatably mount the workpiece in the lathe.\nConventionally, the hole is drilled by mounting the workpiece in the lathe and using a Jacob's chuck and a centre drill. However, this is a relatively complicated centering procedure since the Jacob's chuck is first installed in the lathe, the centre hole is then drilled in the workpiece, the Jacob's chuck is then removed, a centre is then installed in the lathe, and finally the centre is engaged with the drilled hole in the workpiece.\nTo simplify the centering procedure at the lathe, some workshops have a drilling apparatus remote from the lathe to drill the centre hole. However, centering the workpiece on the lathe takes a relatively long time using such drilling apparatus.\nIt is an object of the invention to provide a centering assembly for a lathe which permits a centre hole to be drilled on the lathe and permits a centre to be located in the drilled hole in an improved, faster and simplified manner.\nMy invention includes a centering assembly for centering a workpiece in a lathe including a rotatable centre having an opening along the axis of rotation thereof, a drill in said opening, and means for moving the drill along said axis of rotation from a retracted position within the centre to an extended position for drilling said workpiece.\nPreferably holding means are provided to prevent rotation of the centre when the drill is moved to the extended position.\nFurther objects and advantages of the invention will appear from the following description of a preferred embodiment thereof given by way of example with reference to the accompanying drawings."}
{"text": "In many situations, items must be lifted onto an elevated surface or lowered from an elevated surface. Often, equipment must be lifted onto or lowered from rooftops or elevated platforms. Unless a hoist is provided on the elevated surface, the equipment is typically carried up a ladder by hand or, if the equipment is sufficiently heavy, is lifted onto the elevated surface by a crane.\nHand carrying equipment via a ladder can be dangerous, especially if the equipment is cumbersome or heavy. The use of a crane to move equipment onto or from an elevated surface is not typically dangerous but can be quite expensive. Further, it is difficult in some circumstances to position a crane within close confines so that it can reach the elevated surface.\nThus, a need exists for a lifting device which can be inexpensively and readily positioned to lift items onto and lower items from elevated surfaces."}
{"text": "Modern inventory systems, such as those in container yards, warehouses, superstores, mail-order, e-commerce warehouses and larger manufacturing facilities, use dedicated crane and gantry storage. Such environments also utilize retrieval systems that require dedicated space and large capital investment or vertical pallet or container stacking served by manually driven fork or container hauling trucks. The former provides fast, accurate responses to requests for load, unloading and inventory of items in transit but is a large integrated capital and space investment. On the other hand, manually driven container vertical lift truck based storage systems or yards can cause delays and backlogs in the process of responding to load, unload and inventory requests, and furthermore require drivers and loaders and their salaries, benefits, and management burden.\nHistorically, most loading, unloading and inventory systems are based around vertical storage stacking using pallets or other standardized containers. This arrangement offers a compromise between easy access and three-dimensional storage to reduce area footprint. This type of storage is accessed by forklifts and container handlers of various configurations. The typical way this type of three-dimensional storage is automated is through gantry style S&R or crane units that provide three-dimensional access to the pallet storage locations along a fixed set of travels (i.e., row selection, perhaps through an automated conveyor; column selection with a pallet transport device moving along a preplaced rail, and vertical pallet/container retrieval or storage along an elevator lift mechanism).\nThe type of automation just described is expensive, requires rework of the entire storage/retrieval space, and is an all or nothing proposition—one cannot practically partially automate a space—it either is automated or remains manually retrieved with forklifts or container handlers.\nThe invention disclosed in our prior U.S. Pat. No. 8,965,561, Jacobus et al. describes an approach that assumes use of conventional forklifts and similar container transport platforms, some of which may be equipped with automation enabling them to operate safely along with manually drive lifts into and out of the warehouse. The '561 Patent discloses manual systems enhanced with automation to support more productive automated load acquisition and placement, as well as to improve operator safety by cueing proximal obstacles to the operator in real time.\nThe '561 Patent, further describes pallet engagement and disengagement as beginning with inventory requests; i.e., requests to place palletized material into storage at a specified lot location or requests to retrieve palletized material from a specified lot. These requests are resolved into missions for autonomous fork trucks, equivalent mobile platforms, or manual fork truck drivers (and their equipment) that are autonomously or manually executed to effect the request. Automated trucks plan their own movements to execute the mission over the warehouse aisles or roadways, sharing this space with manually driven trucks. Automated units drive to planned speed limits, manage their loads (stability control), stop, go, and merge at intersections according human driving rules. The automated units also use on-board sensors to identify static and dynamic obstacles and people and either avoid them or stop until potential collision risk is removed. Safety enhance manual trucks use the same sensors to provide the operator with collision alerts and optionally automate load localization and acquisition behaviors for productivity enhancement.\nAutomated or partially automated trucks have the ability to execute task specific behaviors as they progress along the mission, including visiting palletizing/depalletizing, refueling, based upon locations in the order defined by the mission. The vehicles have the ability to drive into or through tight openings using sensor-based driving guidance algorithms, and perform pallet finding and engagement and disengagement onto pallet stacks and into shelving units. Load control automated trucks also can identify pallet locations and pallet identification codes by standard commercial means including but not limited to RFID, barcode, UID, optical characters on printed labels through computer vision/camera, and RF code readers. Each automated truck can be placed into manual drive mode at anytime, and along with unmodified manual fork trucks can be driven over the shared space concurrently used by automated trucks. The automated lift can sense pallets or load positions to control the loading and unloading (or stacking/unstacking) operation, can read pallet identifications (barcodes, RFID, printed labels), can identify moving and fixed obstacles, and can plan from-to paths through using data in the warehouse map.\nManually driven trucks can also be enhanced to include obstacle proximity sensing systems that, instead of directly changing the truck path behavior, alert the operator as to presence or unsafe proximity to the obstacle, indicating a collision risk and even partial automation that might change the trucks speed or cause it to initiate automated braking. Similarly, partial or complete automation of the load acquisition or placement function of a manual truck can be included that operate as indicated for fully automated trucks in the previous paragraph to support more automated and efficient load pick-up and placement operations.\nThe actual load pick-up and placement operations (palletizing and depalletizing) are described as utilizing two front mounted laser sensors are shown in FIG. 1 (FIG. 6 in U.S. Pat. No. 8,965,561). One sensor is mounted on the mast near the floor (34). This unit can be scanned up and down at a small angle by tilting the mast forward and reverse (35). This sensor is primarily tasked with detecting forward driving obstacles, but can also be used to find pallets and pallet stacks. A second sensor is fixed to the fork elevator (36) and sweeps a half cylindrical area (37) as the lift moves up and down the mast. This sensor is tasked primarily for detecting pallet fork openings and the top pallet in a pallet stack, but it too can aid system in sensing obstacles when driving forward. It is noted that between these two sensors the automated truck has almost 100% unobstructed obstacle visibility and is thus, far superior to a human driver's visibility because he/she has to sight through the front forks.\nFIG. 2 shows placement of several video based sensors used to find and read a pallet identifier. The top camera (39) images the entire side of the pallet (40) and uses image processing to identify the gross placement of the bar code, UID, Character notations (OCR) or other identification markers. Then the two smaller field cameras (41) are used to hone in and read the code identifier at sufficient resolution to achieve reliable code deciphering. Alternatively, or in addition, these sensors could include a short range RFID reader or any alternative pallet identification method. An advantage of using cameras is that they can also be used to place roadway or path segment signage up in the warehouse.\nFIG. 3 shows a two-dimensional type of barcode (42) that is in the public domain [Artool Kit, 2012] that can code both content data (the internal bit pattern codes the content) and location of the lift truck relative to the sign placement location. For this code a camera (43) and simple image processing algorithm finds the black square frame, resolving the location of its corners. It then reads the internal bit pattern adjusted for perspective skew. If the deciphered code is “correct”, i.e. can be validated as having been assigned to a known location, its location when added to the offset from the code to the lift truck provides the truck automation controller an exact fix within the warehouse area. These codes could also be placed on the floor as easily as anywhere other known location in the warehouse.\nIn FIG. 4, summarizing automated truck functions, pallet engagement (73) and disengagement (74) as described use range data sets, but in an engagement task specific manner. Pallet engagement is performed in three steps. First the lift truck approaches the approximate position of the pallet stack provided to it as the source or destination location (with orientation) defined by the system controller and provided to the truck in the mission definition. In the second step the truck stops and evaluates its local obstacle map (FIG. 10) for the nearest pallet stack at or beyond its stop position, which it then aligns to. In the third step, the truck uses its mast-controls-by-wire (75) and sensors to:                1. Scan and capture the pallet stack from bottom to top (or to a predefined maximum height)        2. Determine the precision position of the pallet stack face        3. Move the forks to the correct position for stack approachment (for picking a pallet this is alignment with the fork holes of the topmost pallet; for placing a pallet this is alignment to the space just above the topmost pallet)        4. If a pick operation, it optionally uses the wide area camera (39) to find the pallet identification barcode and one of the narrow angle cameras (41) to identify the pallet by reading its pallet identifier. If placing a pallet, lot identification is done on any pallet that is down from the topmost pallet in the stack. Alternative means of identification can be use with or instead of these methods, including RFID, etc.        5. Upon pallet or lot identification, the lift moves forward to align the payload pallet to the stack (or to insert the forks into the topmost pallet holes)        6. To engage, the lift truck lifts the topmost pallet up and pulls back away from the stack; to disengage, it lowers the forks until the pallet rests on the top of the stack and then the lift slightly lifts up and pulls back away from the stack.        7. It then lowers it forks to the predetermined safe travel level above the floor, plans its next movement, and executes the plan, to take the lift and its payload to the next destination (or back to the source location).        "}
{"text": "1. Field of the Invention\nThe field of this invention relates to repairing holes and rusted areas of automobile and truck body panels.\n2. Prior Art\nA major problem concerning the typical automobile whose body is fashioned from sheet metal is that certain sections of the sheet metal body tend to rust away. This rusting causes the automobile to become unsightly, which can lead to the owner's disposing of the vehicle long before the mechanical parts wear out.\nCoping with this problem has long been a challenge to those wishing to repair the automobile body.\nThe method of repairing automobile bodies has changed quite materially over a period of years. Some years ago, it was common practice to repair a dented surface or a hole in the body by flattening out the surface as well as possible by hand, and then applying molten lead to the metal surfaces after the paint had been removed. This molten lead adhered to the metal surface and was built up to a thickness sufficient to produce a smooth outer surface. The required area was sanded to provide a smooth continuous surface and the repaired area was then painted to the color of the original body.\nThis process of \"leading\" the injured areas, was not only time consuming, but it also required a relatively high degree of skill to accomplish good results. As a result, this method has fallen from general use and has been replaced by alternative methods.\nThe current methods of repairing rust damage fall into two categories which can be characterized as \"temporary\" and \"semi-permanent.\"\nThe temporary method involves applying a pressure sensitive aluminum or stainless steel tape or a fiberglass patch over the damaged area, filling over this patch with a polyester filler, priming the area and painting it. This type of repair, while quick and inexpensive, will generally last only about six (6) months. Such a repair, therefore, provides only a temporary solution to the problem which is generally less than satisfactory.\nThe semi-permanent method involves fashioning a metal patch to cover the area in question, welding or riveting the patch in place, (or in some severe cases of rust, replacing the entire section with a new section), filling over the area with polyester filler, priming and painting.\nThis method does not generally provide a fully permanent repair since it does not correct the situation which caused the rust in the first place.\nRusting in an automobile is usually caused by dirt, salt and moisture becoming trapped in a crevice or boxed-in section of the automobile body, so that the body metal is kept constantly damp and susceptible to oxidation. Thus, the rusting process occurs from the backside of the body panels and it is obvious that patches on the surface of a panel do nothing to impede the resuting process and, therefore, such patches are classified as semi-permanent.\nThis method is less than satisfactory not only because it is semi-permanent, but also because it presents the following problems.\n1. A fair amount of skill is required to fashion the metal patch so that it conforms to the original contours of the automobile body. In some instances such as the dog-leg area (the lower area where the rear fender meets the door) the compound curves are so complex that even skilled craftsmen have difficulty in consummating repairs.\n2. Skill is also required in attaching the patch or replacement section to the automobile. The repair may be riveted or welded and care must be taken that the original contours are followed. Welding will most times cause metal to distort.\n3. In most cases, the patch metal is fastened to the outside of the old metal. Thus a ridge, equal to the thickness of the patch material, is created. Since this ridge must be blended out so that the repair will be pleasing to the eye, the polyester filler must extend a relatively large distance from the repair. This necessitates that the paint be ground off a large area around the repair so that the polyester fill can be graded up to cover the patch and approximate the original contour.\n4. Welding or riveting equipment is required.\nIn many instances, repair of rusted areas is so time consuming, and the cost of skilled craftsmen is so high that the labor cost will result in a price so high that such a repair is impractical.\nClearly, a method of repairing holes and rusted areas of automobile and truck body panels which is permanent, and requires less skilled labor is greatly needed."}
{"text": "Spinal pathologies and disorders such as degenerative disc disease, disc herniation, spondylolisthesis, stenosis, osteoporosis, tumor, scoliosis and other curvature abnormalities, kyphosis and fracture may result from factors including trauma, disease and degenerative conditions caused by injury and aging. Spinal disorders typically result in symptoms including deformity, pain, nerve damage, and partial or complete loss of mobility.\nNon-surgical treatments, such as medication, rehabilitation and exercise can be effective, however, may fail to relieve the symptoms associated with these disorders. Surgical treatment of these spinal disorders includes fusion, fixation, correction, discectomy, microdiscectomy, corpectomy, decompression, laminectomy, laminotomy, foraminotomy, facetectomy and implantable prosthetics. As part of these surgical treatments, spinal constructs, such as, for example, bone fasteners, spinal rods and interbody devices can be used to provide stability to a treated region. For example, during surgical treatment, surgical instruments can be used to deliver components of the spinal constructs to the surgical site for fixation with bone to immobilize a joint. Certain spinal surgery approaches utilize a direct lateral approach to access intervertebral spaces, however, these techniques present certain challenges due to the location of musculature and neural structures embedded therein. This disclosure describes an improvement over these prior technologies."}
{"text": "The invention relates generally to a compound and method for repairing damaged plastic and vinyl covered vehicle dashboards and, more specifically, to repair compounds and a two-step repair method for repairing damaged dashboards covered with plastic, vinyl and other similar materials.\nPlastic materials such as vinyl have been widely accepted for use as a covering for automobile dashboards. Vinyl and similar plastic materials are advantageous in this application because they are impervious to water, they remain flexible over a broad range of temperatures and they can be produced in a wide variety of colors and finishes.\nThe use of vinyl materials and the like as a cover for dashboards in automobiles and other vehicles is substantial. Since in such service it will generally be used for an extended period of time, the likelihood of eventual damage is great. Objects such as automobile keys, pens, pencils and toys or heated objects such as matches or cigarettes are obvious causes of such damage. In spite of the fact that such damage may be confined to a small area, the area of the damage is bound to be in an exposed, highly visible region of the vehicle. In spite of the high visibility of such damage, the expense of replacing the entire plastic panel containing the damaged dashboard discourages most consumers from this type of repair. Such damage thus goes unrepaired serving as the situs for tearing and further damage of the dashboard and lowering the overall value of the car or other vehicle.\nI have devoted my life to developing materials and methods for effecting repairs on plastic and vinyl articles. For example, in 1973 I obtained U.S. Pat. No. 3,713,926 for a basic process for the in-situ repair of vinyl material. This process comprehends the filling of a damaged area with a heat curable material of 70 percent polyvinyl acetate and 30 percent plasticizer. A graining paper is placed over the patch and sufficient heat is applied to cure the repair compound. My invention disclosed and claimed in U.S. Pat. No. 3,887,413 teaches another, more sophisticated method of vinyl repair which includes the use of a protective paste which is dispersed on the surface of the undamaged material adjacent the damaged area. U.S. Pat. No. 3,975,558 also utilizes a heat insulating paste. A repair method utilizing a heat transfer tool is described in my U.S. Pat. No. 4,181,547. From the foregoing, it is apparent that I have expended substantial effort with relation to materials and processes directed to the repair of vinyl, leather and similar materials.\nThe field of vinyl repair is an expanding and increasingly sophisticated industry as evidenced by the patent activity of others. U.S. Pat. No. 4,013,495 relates to a method for repairing pigmented and textured flexible material wherein a plurality of thin layers of liquid vinyl repair compound are serially applied and heat cured. U.S. Pat. No. 4,028,160 describes another method wherein precured plastic material is cut into edge abutting patches and used thereto. U.S. Pat. No. 4,086,113 discloses the use of precured patch material, as disclosed in the previous patent, in a preformed patch which is larger than the area of the damaged portion such that it overlaps the damage. The patches are fused to the surrounding undamaged material.\nIn spite of the innovation of myself and others, certain difficulties still exist. One of the most formidable areas of difficulty relates to the attainment of satisfactory curing and bonding of the repair material. This difficulty is the manifestation of a general problem relating to the proper application of heat. In the processes I have developed, it is frequently difficult to provide an appropriate quantity of heat to properly cure all of the vinyl repair compound in a large and thick area of damage. If the vinyl repair compound has not properly cured, due to insufficient heat application, the service life of the patch will be reduced since the material will not achieve appropriate toughness and the bonding achieved with the adjacent material will be less than optimum which may result in separation of the patch material from the surrounding undamaged material. Excessive application of heat, while ensuring complete curing and bonding of the patch material, will invariably destroy the grain texture of the surrounding material or distort it. Clearly, those processes utilizing precured patch material are intended to alleviate the problems attendant insufficiently cured repair material. However, heat application then becomes an even more exacting process since insufficient heat will result in an insufficient bond between the patch and the surrounding material. Generous heat application such as may be frequently needed to achieve melting and bonding of the patch to the surrounding material may, unfortunately, result in damage to the pattern thereof and distortion of the material just as in processes where heat curing repair materials are utilized.\nBeyond the foregoing difficulties, dashboard repair presents a problem relative to matching the hardness of the repair compound, after curing, to that of the surrounding, undamaged dashboard material. The latter is harder than plastic materials, such as vinyl, which are utilized in so many applications including upholstery, luggage, purses and footwear, for which most repair compounds and methods have been developed. Accordingly, currently available vinyl repair compounds are not suited for dashboard repair because the resultant repair patch is significantly softer than the cover material on dashboards."}
{"text": "The present invention relates to a magnetron for microwave oven.\nAn ordinary magnetron for microwave oven oscillating radio waves of 2450 MHz band is provided with an anode cylinder and a plurality of vanes. The vanes are radially arranged inside the anode cylinder. The vanes are alternately connected together in the direction of circumference by a pair of large and small strap rings soldered with the top and bottom ends of the vanes. A spiral cathode is located along the axis of the anode cylinder in the electronic interaction space surrounded by the free ends of the vanes. Both ends of the spiral cathode are fixed to the output side end hat and the input side end hat respectively. Furthermore, both ends of the anode cylinder are fixed to pole pieces formed like funnels of output and input side respectively.\nIn view of miniaturization of heating appliances per se, resource saving and cost reduction, it is required to miniaturize the magnetron for microwave oven. However, merely miniaturization would cause degradation of various characteristics of the magnetron.\nFor example, if the height of the vane in the direction of the axis and the length of the input stem portion are decreased, cathode back heat by electrons or deterioration of load stability may occur. If the distance between both pole pieces is merely narrowed for the purpose of effective utilization of the magnet in order to reduce the height of the magnetron, electromagnetic coupling between the pole piece and the strap ring is increased and result in temperature rising of the cathode due to increase of cathode back heat by electrons. On the other hand, if the height of the vanes in the direction of the axis is shortened to assure the distance between the pole piece and the strap ring to some extent, load stability may be deteriorated. In addition, if the input stem portion is shortened, cathode back heat by electrons is extremely increased and temperature of the cathode is raised. In extreme cases, a part of the cathode may melt. Therefore, a magnetron which height of the vanes in the direction of the axis is 9 to 10 mm has ever been estimated suitable for miniaturization and favorable characteristics.\nJapanese Patent Application Publication No. 1993-035531 gives attention especially to the distribution of magnetic field in the interaction space and discloses a magnetron that has smaller height of vanes in the direction of the axis by improving the shape, dimensions and electromagnetic coupling with the strap of the pole piece. According to this magnetron, change of the shape and the dimensions of pole piece makes that strength difference in the direction of the axis at the inner end surface of the vanes is less than a predetermined ratio, and equalizes distribution of the magnetic field in the interaction space. Furthermore, by digging the strap ring into inside the groove from the side end of the vane to decrease the electromagnetic coupling with the pole piece, cathode back heat by electrons and deterioration of the load stability rarely take place even if the height of the vanes in the direction of the axis or the input stem portion is diminished.\nHowever, the magnetron disclosed in Japanese Patent Application Publication No. 1993-035531 has a vane of more than of equal to 8.5 mm in height in the direction of the axis. If the height of the vane in the direction of the axis is further decreased, it is conceivable that the load stability will be extremely deteriorated, so that the magnetron will not be suitable for practical use. Additionally, if the height of the vane in the direction of the axis is diminished, electrons leaking from the interaction space increase (i.e. dark current increases) because the gap between the end hats at the both ends of the cathode and the ends of the vane in the direction of the axis becomes large. Thus, there is a possibility of deterioration of the output efficiency and melting of the pole piece, etc.\nThe distance between the end hats at the both ends of the cathode, i.e. the effective length of the filament must be shortened in order to avoid deterioration of the output efficiency and melting of the vane, etc. However, if the effective length of the filament is shortened, the load stability will deteriorate further and the cathode back heat by electrons will increase."}
{"text": "This application relates to and incorporates herein by reference Japanese Patent Application No. 11-264156 filed on Sept. 17, 1999.\nThe present invention relates to an organic electroluminescent (EL) display device having an organic EL display panel in which a plurality of organic EL elements are provided.\nOrganic EL display devices are widely used recently, particularly in the field of a compact-sized display. The organic EL device is advantageous in that it requires no back-light, has a quick response characteristics, provides a wide viewing angle, and the like. It also provides a higher luminance and requires a lower drive voltage than an inorganic EL display device.\nOrganic EL elements of the organic EL device degrade with time much faster than inorganic EL elements. The EL elements are driven differently from one another with respect to the number of elements being driven at a particular time. The luminance of each EL element will decrease as it is driven more frequently. Thus, the organic EL display device will have some EL elements that are degraded more than others as the use of the organic EL display device proceeds. Thus, the organic EL display device will resultantly have variations in luminance.\nIt is therefore an object of the present invention to provide an organic EL display device that has less variations in the luminance of display.\nAccording to the present invention, an EL display device has a plurality of organic EL elements and a drive circuit which applies a drive voltage to EL elements to be driven for a required image display. The drive circuit is so constructed to equalize a degree of degradation in luminance among the EL elements.\nIn one aspect of the present invention, the circuit periodically applies a dummy voltage to the EL element that is not driven for the required image display operation, while applying the drive voltage to the EL elements that are driven. The dummy voltage promotes degradation of the non-driven EL elements. Preferably, the dummy voltage has a wave form different from that of the drive voltage and a period shorter than that of the drive voltage.\nIn another aspect of the present invention, the circuit modifies the drive voltage applied to the EL elements to be driven. For instance, the circuit may apply an operation characteristics recovery voltage following the drive voltage to the EL element to be driven that is more degraded than the other. The circuit may apply an initial drive voltage to the EL element to be driven that is less degraded than the other, only when a power supply is started for starting the required image display operation. Preferably, the degree of degradation of each EL element is determined base on an accumulated drive period, the number of drives or the like.\nIn a further aspect of the present invention, the circuit applies the drive voltage for a screen saver operation to all the EL elements or only EL elements that are less degraded."}
{"text": "The currently available hand-held mobile devices have higher and higher performance gained by using a multi-core central processing unit (CPU) having two, four or even more cores. However, with the multi-core CPU that has largely increased operating speed, the heat produced by the hand-held mobile device during operation thereof also disadvantageously increases, which necessitates more efficient heat dissipation solution for the hand-held mobile device.\nAccording to a conventional technical means, a piece of graphite material is correspondingly attached to a heat source of the hand-held mobile device to enable good transfer of heat from the device to an external environment via the graphite material. However, to meet the demand for compactness and light weight, the current hand-held mobile devices all have a very limited internal space, which is too small to accommodate an additional graphite material. In other words, the use of the graphite material would inevitably increase the overall thickness, volume or weight of the mobile devices. It is therefore an important and challenging issue as how to provide a heat-dissipation structure that can be mounted in a hand-held mobile device to effectively solve its problem of heat dissipation without increasing the overall thickness, volume or weight thereof."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for manufacturing nuclear fuel rods by introducing elongated fuel pellets into cladding tubes.\n2. Description of the Prior Art\nThe fuel for nuclear reactors generally is in the form of cylindrical sintered uranium pellets. Groups of such pellets are disposed in elongated cladding tubes designed to prevent coolant from making direct contact with the radioactive fuel. This is to prevent the scattering of radioactive material with the coolant. A particular fuel rod may contain up to 500 individual pellets each having a diameter of about 10 millimeters and a length of about 11 millimeters. The cladding tube itself will thus often have a length of up to about 5000 millimeters. In such fuel rods it is desirable to have a very small clearance in the order of a few hundredths of a millimeter between the outer diameter of the pellets and the inner diameter of the tube.\nThe introduction of the pellets into the cladding tubes has in the past been very difficult and time consuming because of the small clearances involved. Ordinarily strings of pellets have been inserted into horizontally positioned tubes and it has been found that if the end surfaces of the pellets are not absolutely perpendicular to the longitudinal axis of the tube, the string of pellets will become disoriented during insertion whereby the friction between the pellets and the inside surfaces of the tube is increased and a cutting phenomena may even result. Moreover, large insertion forces have often been required in the past due to the high density of the pellets."}
{"text": "Advances in micro-miniaturization within the semiconductor industry in recent years have enabled biotechnologists to begin packing their traditionally bulky sensing tools into smaller and smaller form factors, onto so-called biochips. It would be desirable to develop techniques for biochips."}
{"text": "Conventionally, numerous methods using electronic mails have been known as simple means of operating of various information processing apparatuses. For example, Japanese Patent Application Laid-Open No. 10-334002 discloses transferring an electronic mail to another network. Japanese Patent Application Laid-Open No. 6-139164 discloses updating an address book. Japanese Patent Application Laid-Open No. 6-342409 discloses changing a schedule. Further, Japanese Patent Application Laid-Open No. 5-233493 discloses management of computer.\nHowever, in the above respective conventional technologies, application of such technologies to an information processing apparatus capable of voice output is not disclosed.\nOn the other hand, it is generally possible to voice-output a transmitted electronic mail on the receiving apparatus. However, the voice output cannot be performed at a sender's intended timing. Rather, the timing of the voice output of the transmitted electronic mail is determined by the receiving apparatus."}
{"text": "Testing VLSI circuit wafers has become over the years a significant part of the cost associated with the manufacturing of semiconductor products. Such testing generally entails inputting data into a logic circuit on the wafer at one or more primary inputs or testing test access point, and then reading the output results at one or more primary output terminals or output test access points.\nAdvances in VLSI designs has introduced radical changes in the field of testing and the diagnose of faults principally because the requirement of shipping defect free products to the customer has necessitated adding additional testing features to the design that have become part and parcel of the overall final product.\nTechniques known as Design for Testability have been specifically introduced to this end. An example of such a technique is a scan design known as Level Sensitive Scan Design (LSSD) described in U.S. Pat. No. 3,761,695 to Eichelberger, which is today widely used throughout the VLSI design community.\nThe introduction of scan design has facilitated testing VLSI chips and integrated circuit (IC) packages by introducing linking latches and the like in scan chains memory devices scan chains, and alternating combinatorial portions of combinatorial logic between pairs of the scan chains.\nThe LSSD architecture incorporates several basic test concepts, such as the scan design. In such a design most of the device's storage elements, e.g., latches or registers, are concatenated in one or more scan chains and can be externally accessible via one or more serial inputs and outputs. Storage elements that are not in this category are usually memory or other special macros that are isolated and tested independently. Furthermore, this design methodology ensures that all logic feedback paths are gated by one or more of these storage elements, thereby simplifying a sequential design into subsets of combinational logic sections.\nWhile scan chains are useful in determining whether the logic circuitry is functioning properly, the scan chains themselves may also be defective. While such defects may be from defective latches in the scan chain, if the latches are robust (designed to ensure their integrity), then defects are primarily in the wiring connecting the latches. Such defects may be open circuits (a clean break in the wiring), short circuits (a wire contacting another wire inadvertently), or stuck-at faults (the wiring touching either ground or voltage). The most problematic wiring defect is a stuck-at fault, since the latch otherwise appears to function properly, albeit with a constant input value.\nBasic design concepts in conjunction with the associated system and scan clocking sequences greatly simply the test generation, testing, and diagnosability of complex logic structures. Every latch can be used as a pseudo Primary Input (PI) and as a pseudo-Primary Output (PO) in addition to the standard PIs and POs to enhance the stimulation and observability of the device being tested or diagnosed. LSSD latches are typically implemented in a L1/L2 configuration where the L1 or master latch has two data ports and may be updated be either a scan clock or a functional clock. The L2 or slave latch has only one clock input and that clock is out-of-phase with both L1 clocks. Scanning is performed using separate A and B clocks.\nOf particular importance to the field of testing is a technique known as Self-Test, whereby a set of reproducible random binary test vectors are generated, thereby allowing creating testing means applicable to any chip, multi-chip modules, up to a system level, without requiring heavy expenditures of CPU time. Two basic components of a LBIST (Logic Built-in Self-Test) structure are a Liner Feedback Shift Register (LFSR) and a Multiple Input Signature Register (MISR). The LFSR serves as a pseudo random pattern generator that provides the stimuli for the logic to be tested, while the MISR is utilized to generate a unique signature representing the responses from the logic. Ideally the signature for each failing device is different from the signature of a good device after a predefined number of test cycles.\nThe configuration of the scan chain in the LBIST test mode is partitioned into several sub-chains of approximately the same length as described in U.S. Pat. No. 5,150,366 to Bardell et al. These chains are loaded and unloaded serially for each LBIST test. Once in LBIST mode, the scan chain is reconfigured into a number of parallel sub-chains. The pseudo random data loaded in parallel into each sub-chain is supplied by the LFSR and used as test stimuli. Similarly, the state of all the latches in the sub-chains is unloaded serially into the MISR forming a signature representing the compressed data.\nEach LBIST test cycle, in addition to the loading and unloading of the sub-chains, requires timed application of system clocks to launch the test vector from these latches through the combinational logic and capture the resulting response in the receiving latches. Since a typical system design may consist of several system clocks and various path delays, the clock test sequence and timing set-up may be applied multiple times with different clock combinations and timings. Typically, this is accomplished by an on-product clock generation (OPCG) function and LBIST control.\nAn LBIST test interval in turn consists of a relatively large number of these load/unload sequences followed by the system clock cycle. At the end of the interval, the MISR contents or signature is unloaded and compared to an expected signature. Several signature intervals may be applied to achieve the desire test coverage.\nThis LBIST methodology is an effective Design for Test (DFT) that can support structural test from the chip level, various package levels, up to the system level. Some of the benefits associated with this approach include relatively low test data volumes, minimal VLSI test system requirements, at-speed test rates, and extendibility to system test.\nThe GPTR chain has been getting longer with each server design version. Currently, some of the chains have 1000+ SRLs. GPTR scan chains are expected to become significantly longer as multiple cores and additional arrays on the chip are added. The scan-only GPTR latches perform multiple functions, such as dynamic configurability, redundancy and repair, timing and performance optimization, test and debug “chicken switches”, and many other future design options.\nAlthough the current GPTR chains represent a relatively small percentage of the total latches on the chip, their impact do diagnosability is more severe. In a stuck-at GPTR chain, the configurability of the chip is questionable and, in most cases, diagnosing most other problems on the chip becomes difficult and time consuming. In cases where the GPTR chain causes a very low or zero yield for the product, the diagnosability becomes a severe problem.\nFor illustrative purposes, and with reference to FIG. 1, there is depicted a prior art GPTR scan chain 100 including a chain of shift Register Latches (SRL), SRL1-SRLN, each including a master latch L1 (102) and a slave latch L2 (104). The master latch L1 (102) has a pair of data ports SCAN and DATA that may be captured by the latch responsive to a scan clock A-CLK. The slave latch L2 (104) captures the value stored in the master latch L1 (102) responsive to a second scan clock B-CLK.\nGenerally, the major drawback of the scan based design test methodology is encountered when a scan chain does not function properly which reduces the access to the internal logic of the device, severely complicating the diagnostic process and inhibiting rapid determination of the problem's root cause. In low or zero yield situations, the most common failure is often caused by the scan chain. Since the design requires a scan based design, it is found that the scan chains occupy a significant portion of the real estate area. Having a solution that speeds the ‘stuck-at’ scan chain diagnostics on failing chips results in timely yield improvements and, consequently, ensuring a successful production of the design.\nThe broken or stuck-at scan chain problem is even more severe when it occurs in scan chains that are used to control the test modes and configuration of the remaining latches on the device. In these situations, the diagnosis of the problem is necessary to identify and rapidly resolve catastrophic zero yield or design problems.\nAdditionally, these type of problems are usually encountered early in the technology's life cycle and diagnosing them is critical in improving the process so it quickly achieves manufacturing yield levels. The inability to improve the technology and yield of the device can greatly impact a program or at least severely minimize the revenue that could be realized. Rapid diagnosis to a location for Physical Failure Analysis (PFA) is needed to understand and correct the process anomalies.\nA common approach to diagnose “broken” or stuck-at scan chains is to use the system port to switch the latches down stream from the “break” and than unload the scan chain and determine the last latch position that switches to the complement state of the stuck-at state. Of course, one needs to provide the proper logic state to the system port input of the latch. Furthermore, to get an accurate diagnostic call for the stuck latch, one needs to provide the proper logic input to the next latch downstream to the stuck-at latch, otherwise only an approximate latch diagnostic location can be determined.\nThus, there is a need for a GPTR and a system capable of diagnosing broken scan chain defects of long non-scannable register chains (GPTR), primarily those exceeding 10K latches)."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device and, more particularly, to a semiconductor device having a structure in which a semiconductor chip is covered with an encapsulating resin layer.\n2. Description of the Related Art\nSemiconductor devices such as a large-scale integrated circuit (LSI) are required to have high speed. To meet this requirement, it is recently being attempted to use low-resistivity materials as wires and decrease the relative dielectric constant of an interlayer insulating film. More specifically, the wiring material is being changed from Al to Cu. Also, as an interlayer insulating film, it is being considered to use a so-called “low-k film”, such as a fluorine-doped SiO2 film or an SiO2 film containing an organic component, which has a relative dielectric constant smaller than that of a dense and pure SiO2 film. Note that “an interlayer insulating film having a relative dielectric constant smaller than that of an SiO2 film” used in this specification and in the scope of claims means “a low-k film”.\nThe relative dielectric constant of an interlayer insulating film can be decreased by making the density of the film lower than the intrinsic density of the material, or by eliminating the dielectric polarization in the material.\nFor example, a low-density insulating film is generally obtained by making the material porous. The relative dielectric constant of an insulating film thus obtained is made smaller than the intrinsic relative dielectric constant of the material by pores in the insulating film. Unfortunately, these porous insulating films generally have low mechanical strength.\nIn addition, an insulating film from which the dielectric polarization is eliminated has low adhesion to another insulating film. That is, in a structure in which a plurality of interlayer insulating films including an interlayer insulating film having a low relative dielectric constant are stacked, peeling readily occurs at the interfaces between the insulating films.\nThe low mechanical strength and low adhesion of an interlayer insulating film pose problems such as film peeling in steps in which heating is performed. Examples are packaging of a semiconductor chip including a multilayer interconnection structure, mounting the obtained semiconductor package on a printed wiring board, and a temperature cycling test (TCT). This is so because in these heating steps, a thermal stress caused by a difference in coefficient of linear expansion between a semiconductor substrate and an encapsulating resin layer is applied to the multilayer interconnection structure to destroy or peel an interlayer insulating film.\nJpn. Pat. Appln. KOKAI Publication No. 9-246464 describes a semiconductor package in which a semiconductor chip is mounted on a wiring substrate via insulating resin by flip chip bonding, and the obtained structure is encapsulated with an encapsulating resin. In this semiconductor package, the insulating resin covers the side surfaces of the semiconductor chip. Jpn. Pat. Appln. KOKAI Publication No. 9-246464 also describes that a resin having Young's modulus smaller than that of the encapsulating resin is used as the insulating resin. The invention described in Jpn. Pat. Appln. KOKAI Publication No. 9-246464 can prevent, by the use of this structure, destruction of connecting portions occurring between the semiconductor chip and wiring substrate owing to, e.g., the thermal stress of the encapsulating resin."}
{"text": "A push for higher performance and smaller die size drives the semiconductor industry to reduce circuit chip area by approximately 50% every two years. The chip area reduction provides an economic benefit for migrating to newer technologies. The 50% chip area reduction is achieved by reducing the feature sizes between 25% and 30%. The reduction in feature size is enabled by improvements in manufacturing equipment and materials. For example, improvement in the lithographic process has enabled smaller feature sizes to be achieved, while improvement in chemical mechanical polishing (CMP) has in-part enabled a higher number of interconnect layers.\nIn the evolution of lithography, as the minimum feature size approached the wavelength of the light source used to expose the feature shapes, unintended interactions occurred between neighboring features. Today minimum feature sizes are approaching 45 nm (nanometers), while the wavelength of the light source used in the photolithography process remains at 193 nm. The difference between the minimum feature size and the wavelength of light used in the photolithography process is defined as the lithographic gap. As the lithographic gap grows, the resolution capability of the lithographic process decreases.\nAn interference pattern occurs as each shape on the mask interacts with the light. The interference patterns from neighboring shapes can create constructive or destructive interference. In the case of constructive interference, unwanted shapes may be inadvertently created. In the case of destructive interference, desired shapes may be inadvertently removed. In either case, a particular shape is printed in a different manner than intended, possibly causing a device failure. Correction methodologies, such as optical proximity correction (OPC), attempt to predict the impact from neighboring shapes and modify the mask such that the printed shape is fabricated as desired. The quality of the light interaction prediction is declining as process geometries shrink and as the light interactions become more complex.\nIn view of the foregoing, a solution is needed for managing lithographic gap issues as technology continues to progress toward smaller semiconductor device features sizes."}
{"text": "Mënière's disease is a disorder of the inner ear. It is characterised by an endolymphatic hydrops of unknown cause. Symptoms consist of disturbed balance, vertigo, impaired hearing and tinnitus. It has been shown that people suffering from Mënière's disease can achieve relief from symptoms of dizziness and nausea through the application of elevated pressure to the outer and middle ear.\nIt is theorised that an air pulse, which applies the pressure to the round window of the inner ear may assist in reversing the flow of toxic potassium ions from the scala tympani to the scala vestibuli. The present invention provides a low cost device designed to allow users to apply elevated pressure levels to their own ears in order to achieve such relief.\nThe inventor has also recognised that causing a negative change of pressure in the ear canal may be useful in relieving pain arising from other causes, such as unequal pressure between the middle ear and outer ear that occurs during descent of an aircraft. The present invention also provides a device which may achieve such relief."}
{"text": "1. Field of the Invention\nThis invention relates to methods and apparatus for systems for manufacture of semiconductor devices and more particularly to yield loss analysis in yield management systems.\n2. Description of Related Art\nU.S. Pat. No. 5,598,341 of Ling et al. for \"Real-Time In-Line Defect Disposition and Yield Forecasting System\" shows a real time defect disposition and yield forecasting system which uses KLA and INSPEC inspection tools. A particle in contact with two conductor lines is described as a \"killer defect.\" It is stated that the particle \"may kill or prevent the normal operation of the semiconductor device which utilizes conductors . . . \" Ling describes a real time, in-line system. An inspection method/tool provide for inspection of at least two layers of a semiconductor wafer, producing \"first information\" on particle size, particle location, and the number of particles. A design review station inspects those at least two layers and produces \"second information\" including the layouts of each of those at least two layers. A yield management station is operatively connected to the in-line fabrication inspection tool and to the design review station. The yield management station determines at least one of a number of killer defects for the devices in each of the at least two layers or a defect sensitive area index for each of the at least two layers responsive to layers using the first and second information. Determination of a priority is made for analyzing each of the at least two layers responsive to at least one of the number of killer defects and the defect sensitive area index for each of the at least two layers. See col. 5, line 55 to col. 7, line 8. Ling does not appear to disclose the 6th order polynomial curve fit to identify the killer stage, but Ling does describe a general analysis method/tool/program, Col. 6, line 55 to Col. 7, line 8 .\nU.S. Pat. No. 5,475,695 of Caywood et al. for \"Automatic Failure System\" teaches an automatic failure analysis system which discloses automatic fault extraction (AFE) and a specific analysis technique called inductive fault analysis (IFA).\nU.S. Pat. No. 4,801,869 of Sprogis for \"Semiconductor Defect Monitor for Diagnosing Processing-Induced Defects\" shows a design for a defect monitor pattern/wafer.\nU.S. Pat. No. 5,544,256 of Brecher for \"Automated Defect Classification System\" shows an automated defect classification system for analyzing digital images of defects.\nKLA In-line defect inspection tools and the Yield Management System (YMS) are popular tools for studying the visual defect in process but heretofore, there has been no systematic method for determining the correlation between yield loss and visual defects. In the more complicated manufacturing processes, the state of the art methods of yield loss analysis have made it difficult to correlate the major killing stages to the resulting losses in yield. Thus it has been difficult to be efficient in selecting which action to take to improve the yield."}
{"text": "Today, we are more connected as a society than ever before. Data is continuously being mined and stored from various sources by a plethora of companies and individuals. Data may be, among others, data from any type of sensor, data tracked by companies or data relevant to the public at large. Examples of data affecting the public at large may be traffic data, weather data, stock price data, etc.\nCompanies often use sensors to track the condition of their equipment and inventory. For example, sensors are used in oil depots (“tank farms”) to monitor various statistics of the oil stored at the tank farms as well as conditions of the equipment storing the oil. Sensors may be used to track the amount of oil stored in each storage tank, the water density in each tank, the inflow into each storage tank, the outflow of oil from each tank, etc. In addition, sensors may be used to track the temperature of the interior of the tank, the exterior or the humidity surrounding the tank.\nIn addition, companies track their inventory and sales at their distribution centers. For example, an oil distribution company will track the amount of oil it sells to each gas station, airport, shipping yard, etc. The company likely tracks the price at which each barrel of oil was sold, the date of the sale, etc. The company may also track its supply chain and distribution processes such that the time and steps taken to refine the oil as well as the location of each transport vessel (e.g., ship or truck).\nFurthermore, some forms of gathered types have been used to predict future events. For example, weather data, e.g., data relevant to the public at large, is routinely collected and used to predict future weather systems in a given geographic area. For example, data may be collected from thermometers, barometers, windsocks, humidity sensors, air pressure sensors, etc.\nHowever, even though all of this data may be collected and stored by various sources, the use of such data in predictive analytics has thus far been limited. For example, the data used to predict the weather forecast (e.g., data relevant to the public at large) has not been combined with data collected by companies regarding their inventory and shipments (e.g., business applications data such as the enterprise resource planning (ERP) of the company)."}
{"text": "The present invention relates to a multi-piece golf ball having at least a three-layer construction composed of a core, an intermediate layer and a cover layer. More specifically, the invention relates to a multi-piece golf ball having a high rebound resilience and durability, and having an excellent feel on impact. The invention also relates to a multi-piece golf ball which has a low spin rate and can travel a good distance even when hit at a low head speed.\nThe golfing population has grown in recent years, and the desires of players with regard to golf balls have become increasingly diverse and personalized. To address such desires, various investigations are being conducted on ball construction. Most players are amateur golfers with head speeds (HS) in a range of generally 30 to 45 m/s. Amateur golfers derive greater enjoyment from playing the game when the ball travels farther on shots with a driver and has a good, soft feel at the moment of impact. In addition, because amateur golfers have a head speed that is somewhat lower than that of professionals, it tends to be harder for them to get enough lift on the ball. Amateurs thus use clubs having a relatively high loft angle, as a result of which balls hit with a number one wood (W#1) tend to take on too much spin and acquire an excessive spin rate.\nJP-A 2005-218858 teaches a golf ball having a core, an intermediate layer that encloses the core, and a cover that encloses the intermediate layer. By having the intermediate layer and the cover each formed of an ionomer resin-containing thermoplastic resin, and by concurrently optimizing the respective values for the core diameter, deflection of the core under 100 kg of loading, intermediate layer thickness, Shore D hardness of intermediate layer, cover thickness, Shore D hardness of cover, deflection of the golf ball under 100 kg of loading, and the difference obtained by subtracting the deflection of the golf ball under 100 kg of loading from the deflection of the core under 100 kg of loading, such a golf ball achieves a combination of a good, high rebound, low spin and good feel when used by the ordinary amateur golfer who values distance and feel more than spin performance.\nJP-A 2004-16583 discloses a golf ball having, in order from the inside: a solid core of at least one layer, an intermediate layer, and a cover. In this ball, a good rebound and a soft feel on impact are achieved by optimizing the respective hardnesses and thicknesses of the core, intermediate layer and cover.\nHowever, these golf balls leave room for further improvement in their distance of travel and feel on impact. In particular, the challenge has been how to lower the spin rate on the ball when hit with a driver in the low head speed range typical of the amateur golfer, and thereby improve the distance traveled by the ball."}
{"text": "A γ′-precipitation strengthened Ni-based superalloy is used as, for example, high temperature parts for a gas turbine or a steam turbine that requires mechanical strength under high temperature environment. It is said that the γ′-phase is composed of Ti, Al, Nb, and Ta and that a precipitation amount thereof can be increased by increasing a content of these constituent elements in the alloy and thereby mechanical strength of the alloy at high temperature can be enhanced.\nOn the other hand, in the case where the precipitation amount of the γ′-phase is made large so as to increase the mechanical strength of the alloy at a high temperature, the hot forgeability (hot workability) of the alloy in the production process decreases and, if deformation resistance is thereby made excessively large, the forging itself cannot be performed in some cases. Particularly, it becomes a large problem in a large-sized product such as a turbine disk in which deformation by hot forging is unavoidable. Accordingly, a component composition of an Ni-based superalloy having both of the high-temperature strength and the hot forgeability has been investigated.\nFor example, Patent Document 1 discloses, as such an Ni-based superalloy, an alloy containing, in terms of % by mass, Al of from 1.3 to 2.8%, Co of from a minute amount to 11%, Cr of from 14 to 17%, Fe of from a minute amount to 12%, Mo of from 2 to 5%, Nb+Ta of from 0.5 to 2.5%, Ti of from 2.5 to 4.5%, W of from 1 to 4%, B of from 0.0030 to 0.030%, C of from a minute amount to 0.1%, and Zr of from 0.01 to 0.06%, in which, in terms of atomic %, (1) Al+Ti+Nb+Ta is from 8 to 11 and (2) (Ti+Nb+Ta)/Al is from 0.7 to 1.3, Therein, it is said that the total amount of Al, Ti, Nb, and Ta defines the solid solution temperature of the γ′ phase and the γ′ phase fraction, and according to the expression (1), the γ′ phase fraction is controlled within a range of from 30 to 44% and the solid solution temperature is controlled to lower than 1145° C. Furthermore, it is said that, according to the expression (2), the mechanical strength under high temperature environment owing to the γ′ phase is enhanced and also the precipitation of harmful η-type and 8-type needle-like intermetallic compound phases is prevented. It is said that according to the above, the alloy has such a high forgeability that cracking is not generated even in the forging at a temperature higher than the solid solution temperature of the γ′ phase, which is impossible in the case of UDIMET 720 (“UDIMET” is a registered trademark), and also said that the mechanical strength at 700° C. that is an operating temperature of a turbine can be increased as compared with the case of the Ni-based superalloy called 718 Plus.\nMoreover, Patent Document 2 discloses an Ni-based superalloy having a component composition containing, in terms of % by mass, C of more than 0.001% and less than 0.100%, Cr of 11.0% or more and less than 19.0%, Co of 0.5% or more and less than 22.0%, Fe of 0.5% or more and less than 10.0%, Si of less than 0.1%, Mo of more than 2.0% and less than 5.0%, W of more than 1.0% and less than 5.0%, Mo+½W of 2.5% or more and less than 5.5%, S of less than 0.010%, Nb of 0.3% or more and less than 2.0%, Al of more than 3.00% and less than 6.50%, Ti of 0.20% or more and less than 2.49%, in which, in terms of atomic %, Ti/Al×10 is 0.2 or more and less than 4.0 and Al+Ti+Nb is 8.5% or more and less than 13.0%. Particularly, in Patent Document 2, the precipitation amount of the γ′ phase is increased by increasing the addition amount of Al, Ti, and Nb and, it is described that the high-temperature strength and the hot forgeability are in a trade-off relationship. In Patent Document 2, it is said that the content of Al is increased to prevent the solid solution temperature of the γ′ phase from rising and the high-temperature strength and the hot forgeability are both achieved.\nPatent Document 1: JP-T-2013-502511\nPatent Document 2: JP-A-2015-129341"}
{"text": "Computer service vendors provide a wide range of services, and an increasing number of companies are becoming consumers of the various services. Services available today include telecom (telecommunications, datacenter network services, mobile/edge communications services), computational services (compute), both application layer and infrastructure layer, and data storage (storage).\nUnfortunately, billing standards and methods differ greatly between different vendors and between different types of services. Different types of services may be billed by flat-rate, by the day, week, or month, by the byte (or multiple thereof), or on a per-minute basis. Other factors include data speed and tiered service levels.\nThese varying billing methods cause unnecessary duplication, overhead, and other management issues.\nThere is, therefore, a need in the art for a system, process and data format for unified billing for computing services."}
{"text": "The widespread availability of police \"scanners\" and other inexpensive consumer radio receivers has created significant security problems for law enforcement and other radio users. For example, it is now possible for criminals to monitor police radio communications in order to discover police whereabouts and activities--and thereby elude the police. Although some police forces have adopted the practice of talking in \"code\" to make their communications less understandable, these codes often make informative conversation more difficult--and the \"codes\" typically are relatively easy to \"decode\" after listening to police conversations for a few weeks.\nModern state-of-the-art mobile radios are digital. They convert the user's voice into a digitized data stream consisting of \"bits\" having \"1\" and \"1\" values before sending the \"bits\" over the radio channel--and similarly they receive communications in digitized data stream form and convert the received \"bits\" into analog voice signals for reproduction by a loud speaker. Although the use of digitized speech transmission prevents most police \"scanners\" from receiving intelligible signals, more advanced \"digital\" radio receivers available to consumers or criminals may still allow eavesdropping. Moreover as the cost of \"digital\" scanning receivers comes down, this type of receiver will become more widespread. In addition, a certain class of radio users (e.g., people associated with the FBI, CIA, military and other highly sensitive organizations) require an extremely high degree of communications security. Thus, there is need for communications that provide greater security and protection against eavesdropping than can be achieved using digitized speech.\nIn response to this need for greater security, major manufacturers of land-mobile radio equipment in the United States have for many years incorporated \"encryption\" into their radio products. The state of radio electronics technology has advanced such that it is now practical and cost-effective to have the radio equipment automatically electronically \"code\" (\"encrypt\") and \"decode\" (\"decrypt\") communications electronically.\nBriefly, \"encryption\" and \"decryption\" are part of \"cryptography,\" the art of communicating securely in the presence of an \"enemy\" or \"attacker\". \"Encryption\" takes a \"clear\" text message and transforms it into unintelligible form (\"cipher text\"). \"Decryption\" reverses the process, transforming the cipher text back into the original \"clear\" text. In modern \"cryptosystems\" such as \"DES\" (\"Data Encryption Standard\"), it is computationally infeasible to derive the \"clear\" text from the cipher text unless one knows beforehand the particular transformation that should be employed. So long as the decryption transformation remains a closely guarded secret, parties to an encrypted communication can feel safe knowing that only authorized people (i.e., the people who know the specific decryption transformation to use) will be able to decrypt the encrypted communications. That is, even if an attacker successfully intercepts an encrypted message, the attacker will be unable to decrypt the message to recover the \"clear\" text.\nBecause of the difficulty of designing and verifying the security of a cryptosystem, it has become commonplace to define encryption and decryption transformations with two components: (a) an algorithm that defines a family of transformations; and (b) a cryptographic \"key\" that specifies a particular one of many (usually a very large number of) transformations within the family. The cryptographic algorithm can be contained in standard, readily available integrated circuit chips and may even be widely published (e.g., the details of the \"DES\" algorithm is published in Federal FIPS Publication No. 46). However, before the algorithm is used by, for example, Steve and Carole to encrypt and decrypt messages, Steve and Carole select and agree upon at least one cryptographic \"key\" they intend to use with the algorithm. Steve and Carole keep this cryptographic key a closely guarded secret. Since the cryptographic key determines the particular encryption/decryption transformation(s), an attacker who knows the basic cryptographic algorithm but doesn't know the key will be unable to decrypt intercepted encrypted communications. Moreover Carole and Tim can agree to use a different cryptographic key for their secure encrypted communication, and Steve intercepting the communication will be unable to decrypt it successfully. To provide increased security against long-term cryptographic attacks, \"traffic analysis,\" etc., parties can agree to periodically change their agreed-upon cryptographic keys--or they may agree to use a different key if they feel the security of their communications has been compromised.\nEricsson-GE Mobile Radio Communications Inc. (\"EGE\") of Lynchburg Virginia has for many years sold products under the trade name VOICE GUARD which include a digitized voice encryption/decryption capability. Motorola Inc. of Schaumberg Illinois has similarly for many years sold \"SECURE NET\" land-mobile radio equipment including digitized voice encryption. Although prior radio products were limited to only a single encryption technique, more recent digital voice radios (i.e., EGE's AEGIS/VG digital radio products) currently use any one of three conventional key-based encryption o algorithms (DES, VGE or VGS) for securely transmitting digitized voice over an RF channel.\nIn using such systems, there arises a practical problem of how to load the secret cryptographic key(s) into all radio transceivers intended to participate in secure communications. In EGE's prior systems, the \"key\" information is loaded into each radio individually using a device called a \"key loader.\" This keyloader device communicates with the radio over a serial data cable, and downloads \"key\" data into the radio for use in defining the particular encryption/decryption transformation to be used. The user may download \"keys\" to the radio at any time by connecting up the key loader to the radio and specifying the new key data. When the radio receives new \"keys\" from the key loader, it must store and retain the \"keys\" so that they can be retrieved every time the radio is subsequently turned on.\nIn prior radio products, the key storage method depended on the type of encryption algorithm in the radio. Prior EGE radio designs using DES encryption transferred the \"key\" data to a special-purpose DES encryption/decryption integrated circuit \"chip\" for storage. This DES \"chip\" is connected to a small battery (e.g., a lithium cell) that continually supplied power to the DES chip even when the rest of the radio is turned off. The solution of storing the \"keys\" in the DES chip is very secure since the radio microcontroller never \"sees\" (and cannot access) the key information once it is loaded, and any attempts to read the key information out of the DES chip will almost certainly result in erasure of the key information before it can be successfully read. However, this arrangement requires the existence of a battery-backed device to maintain the \"keys\" throughout the power cycle. Moreover, commercially-available DES chips have a limited key storage capacity whereas in the context of a mobile radio communications system it may be necessary for a given mobile radio to select between a number of different keys corresponding to, for example, a number of different secure communications recipients. In prior VGE encrypted radios sold by EGE, the key information was simply stored in a table in EEPROM (\"Electrically Erasable Programmable Read Only Memory\"). This storage solution removed the requirement for a battery-backed device, but was not very secure since the key information is simply in a table in EEPROM and thus can be read out by someone willing to take the time to \"dump\" the contents of EEPROM--a relatively simply process that can be performed using readily available, relatively inexpensive equipment. Because of this \"reverse engineering\" possibility, the security of the entire radio communications system becomes compromised if even a single radio transceiver falls into the wrong hands. Of course, it is always possible in such a system to change encryption keys on a system-wide basis--but the logistical difficulty of reprogramming each individual radio transceiver at a service depot using a key loader would cause secure communications to be disrupted for hours, days or even longer.\nThere have been prior attempts to securely store encryption keys within a mobile radio transceiver. See, for example, U.S. Pat. No. 5,150,412 to Maru, which discloses a mobile radio telephone including a single chip microcomputer (security module) including an internal non-volatile EEPROM encryption/decryption key store. Whenever external access of the EEPROM key store is attempted (e.g., in order to test the function of the key store), circuitry automatically clears the EEPROM contents--thereby preserving the secrecy of the encryption/decryption keys. This technique has the drawback of requiring a specially designed security module with special-purpose circuitry for destroying key store contents when external access is attempted.\nIt would be highly desirable to provide an arrangement for safeguarding the secrecy of encryption/decryption keys stored by a mobile/portable radio transceiver that does not require any additional hardware components or other costly additions to the transceiver architecture and yet provides flexibility in securely storing a large number of selectable different cryptographic keys.\nThe present invention provides a digital radio having a table in non-volatile memory such as EEPROM for \"key storage\" as in prior VGE products, but the \"keys\" are stored in an \"encrypted\" form such that their identities are not readily revealed by a \"dump\" of memory contents. Additional security is provided in accordance with the present invention by extracting the \"keys\" from the stored table and re-\"encrypting\" the entire table each time a key loader device is attached to the radio. This re-encryption adds another level of complexity to the process should someone attempt to \"break\" the cryptosystem by repetitively loading different \"keys\" into the radio.\nIn somewhat more detail, a digital radio provided by the present invention \"hides\" or \"shrouds\" its key store information by in some sense \"encrypting\" the key information before storing it in the radio's internal EEPROM memory. A pseudo-random function is used as part of the shrouding technique. This use of a pseudo-randomization factor means that the keys are shrouded differently from one radio to the next--and that the same radio shrouds the keys differently from one shrouding operation to the next. An unauthorized person trying to gain access to the keys will presumably \"dump\" the entire contents of the EEPROM including the key store, but this information will be useless unless she also knows the particular shrouding transformation used. To learn the shrouding technique, the attacker would have to dump the entire program store and reverse engineer the control program software in detail--an extremely expensive and time-consuming process exposing the attacker to copyright infringement liability.\nA mobile or portable digital radio provided by the present invention first constructs and writes a table containing pseudo-random data into a key store section of internal EEPROM. The radio's cryptographic keys are written \"over\" the random data, and are written at locations which can change from one key loading operation to the next--thereby \"hiding\" the cryptographic keys by \"burying\" them somewhere in a \"sea\" of pseudo-random data. As an additional protection in the preferred embodiment, the keys before being stored are first \"encrypted\" in a sense that they are transformed based at least in part on random data stored in other parts of the table. As a result, the stored cryptographic keys are hidden among a series of random data values, and the stored key data itself \"looks\" like random data. As a result, an attacker would be unable to learn the identity of the stored cryptographic keys from a dump of the EEPROM table unless she knew where to find the stored keys within the table, and unless she also knew what particular transformation should be used to decrypt and thereby recover the keys.\nIn accordance with a further feature provided by the present invention, the entire table randomization and transformation process is repeated every time a key loader device is connected to the radio transceiver. The actual key data is extracted by performing an inverse transformation, and a random number generator is used to re-randomize the table. The key data is transformed using the newly randomized table and the cryptographic keys along with the associated (new) index are redeposited into the table. This means that the keys typically end up being stored in a different place within the table, and that a different decryption/extraction transformation based on the information stored in the table must be used to recover them.\nIn accordance with a further feature provided by the present invention, multiple key banks are used to provide enhanced voice security by increasing the number of encryption keys available for use by a radio. This feature provides the additional advantages that the number of times a radio must be key loaded is reduced, and the number of personality configurations for groups, channels, and systems is greatly increased.\nPrior portable or mobile two-way radios store only a limited number of encryption keys (e.g., seven encryption keys for EGE's Voice-Guard private voice operation, and six encryption keys for EGE's AEGIS private voice operation). Different keys can be selected for channels, groups and special calls. The number of keys to choose from is very limited. Also, if the user feels private voice calls are no longer secure using the programmed key, the only choice the user has is to discontinue communications until the radio can be key loaded with new encryption keys. Key loading the radio can be very time consuming because each radio must be individually connected to the key loader.\nThe preferred embodiment provided by the present invention solves this problem by using multiple banks of encryption keys--all stored in the same random-data EEPROM table described above. The radio can store multiple banks of keys with n (e.g., six or seven) keys per bank to maintain compatibility with existing radios. The key bank to use can be specified on a per system basis using a radio personality. The radio personality can contain the same system data repeated multiple times with only the key bank changing. The following illustrates an exemplary radio personality:\n______________________________________ System Key Key Bank Group Key ______________________________________ SYS1 3 1 Fire 1 SYS2 3 2 Fire 1 ______________________________________ (with key banks 1 and 2 containing different sets of cryptographic keys).\nIf users feel trunked calls on the group \"fire\" are no longer secure on system \"SYS1,\" the users can instantly change to the other system \"SYS2\" and continue encrypted communications using different encryption keys stored in key bank 2. Increasing the number of keys available to the user provides for more configurations on the radio personality. For example, conventional (non-trunked) operation could use the first four banks of keys, and trunked operation could use the second four banks. Different banks and/or keys can be used for different cryptographic modes (e.g., VGS, VGE or DES). In addition, increasing the number of encryption keys that can be stored in the radio can reduce how often the radio needs to be key loaded."}
{"text": "1. Technical Field\nThe present disclosure relates to bases, and particularly to a base for an electronic device and an electronic apparatus using the base.\n2. Description of Related Art\nElectronic devices, such as a mobile phone, a media player, and a personal digital assistant (PDA) generally include video playing and data transmitting functions. However, users need to hold such electronic devices to view the video. Furthermore, the electronic device may need a data cable detachably coupled to the electronic device to receive or transmit data from or to another electronic device. So, it is inconvenient for the users to use the electronic device.\nTherefore, there is room for improvement within the art."}
{"text": "Within each day exists a cycle that affects attention, alertness, energy and mood of a human being. This internal rhythm is produced and regulated by a neurological mechanism called the circadian clock. In healthy humans, the secretion of cortisol from the adrenal glands follows a diurnal cycle with a profound increase after awakening. This increase after awakening, a phenomenon termed the cortisol awakening response, appears to be a distinct feature of the hypothalamus-pituitary-adrenal axis, superimposing the circadian rhythmicity of cortisol secretion. Another circadian effect is the morning surge of blood pressure. A sudden activation of the sympathetic nervous system is the primary mediator of the morning surge. Increased-mediated sympathetic vasoconstriction has been found in normal subjects. Whereas arousal from sleep is associated with a slight rise in plasma epinephrine, arising induces a significant rise both in epinephrine and norepinephrine.\nIn Panza et al., “Circadian variation in vascular tone and its relation to α-sympathetic vasoconstrictor activity”, N Engl J Med 1991; 325:986-990, it is described that there is a circadian rhythm in basal vascular tone, due either partly or entirely to increased alpha-sympathetic vasoconstrictor activity during the morning. This variation may contribute to higher blood pressure and the increased incidence of cardiovascular events at this time of day.\nCurrent methods for circadian phase detection usually require obtrusive sampling methods. One the one hand, it is possible to derive information on the circadian phase from the core body temperature. The core body temperature can be measured with an intra body sensor, such as an e-pill. Having to swallow an e-pill and to retrieve it from excrements is cumbersome. Alternatively, the temperature in the mouth or in the armpit can be measured, which may also be inconvenient to the individual. On the other hand, it is possible to derive information on the circadian phase from a salivary melatonin test. The salivary testing of melatonin requires the regular sampling of saliva. This is a burden to the individual.\nIn DE 4221526 A1 a system for controlling, measuring and training psycho-emotional adaption processes is disclosed. The system serves for non-invasively verifying different psycho-emotional states like arousal, stress, fear, relaxation, impact of ataractics etc. The system evaluates the secretion of the perspiratory glands by means of pulsating direct current. The system is implemented in the form of a wristwatch and converts electric values of the skin conductance into analog angle readings. A high arousal (low skin conductance) corresponds to a high angle reading (high impulse frequency).\nIn Till Roenneberg: “Chronobiology: The human sleep project”, Nature, June 2013, a research activity at the University of Munich is presented. Researchers have gathered a database including entries for more than 150.000 individuals from all over the world that have provided the times at which they go to bed, prepare for sleep, fall asleep, wake-up and get up.\nConsequently, it is desirable to provide an approach for unobtrusive circadian phase detection. Such an approach could be used to provide a gentle adaption of the circadian phase of a person to an external requirement."}
{"text": "The invention relates to a ionizer, device which neutralizes static electricity buildup on synthetic fibers, sheets, webs, when such materials are processed. Generally, an ionizer has the appearance of a bar, placed across movement direction of the charged material, on the side of the surface where the static charge is buildup. The distance between the ionizer and material varies from one to five inches. Electrically powered ionizers are energized by a high voltage power supply, by means of a high voltage cable. Hot ionizers are named that ones which have the high voltage needles directly connected (jabbed) into the conductor of the high voltage cable. A person who accidentally touches a needle can receive an electric shock. The severity of this shock depends on the current rating (usually limited to 5 milliamperes) of the high voltage power supply. Shockless ionizers are named that ones which have high voltage needles capacitively coupled to the high voltage cable. Due to the small capacitance between the high voltage needle and the conductor of the high voltage cable, the current is limited to few microamperes, and a direct touch of a high voltage needle is painless.\nIt is known an ionizer, which, in order to discharge or neutralize static electricity, consists of a row of high voltage needles inserted into the conductor of a high voltage cable, and one or two electrodes connected to the ground, which flank the needles row. A high voltage power supply, of 5 to 15 kV, energizes high voltage needles by a high voltage cable. The other terminal of the high voltage power supply is connected to the ground. Between high voltage needles and grounded electrode/s of ionizer occurs a corona discharge, generating nearby positive and negative ions. The charged material passing through ionized area attracts the necessary quantity of ions of opposite polarity, until its charge gets neutral. This type of ionizer is efficient for materials with a slight or average charge. For the materials with a high or a high large static charge this ionizer is not efficient, because all the ions generated in corona discharge are attracted by the charge of material, but the charge is still not entirely balanced this way. The density of ions is maximum in the area where the distance between high voltage needles and grounded electrode/s is minimum. This fact decreases the efficiency of neutralization.\nIt is also known an ionizer, named induction or passive ionizer (neutralizer), which consists of a row of needles or thin wires connected to the ground. The material moves nearby the ionizer, and his charge creates an electrical field between grounded needles and itself. A breakdown of the air appears when the electrical field at the peaks of the needles exceeds a threshold value. A small electrical discharge occurs between the charged surface and the grounded needles, this way the material gets discharged. The breakdown of the air and small electrical discharge stop in the moment when the value of the electrical field between the charged surface and grounded needles drops under threshold value. Therefore on the material remains a residual charge which is no more possible to be neutralized this way. The induction ionizer is inefficient for the materials with slight static charge, which can not create an electrical field higher than threshold value.\nThere is also known an ionizer which consists of a row of high voltage needles, connected to a high voltage power supply. On both sides, along of high voltage needles row, are arranged two saw-teeth shape electrodes, connected to the ground. Between the high voltage needles and grounded electrodes takes place a corona discharge, generating positive and negative ions. The materials with a slight static charge get discharged by attracting ions of opposite polarity. High static charges get neutralized mainly by small electrical discharges between the peaks of the teeth of the grounded electrodes and charged material. The materials having an average static charge get neutralized by small electrical discharges to the peaks of grounded electrodes, as well by attracting ions of opposite polarity.\nThis type of ionizer is not enough efficient because the density of ions is not uniform. The corona discharge is perpendicular to the axis of the ionizer and occurs mainly where the distance between the high voltage needles and grounded electrodes is minimum. The density of ions is maximum in the area where the corona discharge takes place, and lowers to a minimum at the half of the distance between the needles. Consequently, the neutralization of the electrical charge of the material is good in the areas where the ions density is maximum, and decreases in the areas where the ions density is poor. The result of the decreased neutralization in these areas is an increasing of the residual electrical charge on the material.\nU.S. Pat. No. 4,216,518 refers to a Capacitively Coupled Static Eliminator with High Voltage Shield. It consists of a grounded metallic U shape housing, which supports high voltage electrode (plate configuration) and surrounding insulation, high voltage needles, capacitively coupled with high voltage electrode and high voltage shield. An AC high voltage power supply energizes high voltage electrode and shield. This ionizer is efficient for neutralization of slight or average static charges. The corona discharge current between high voltage needle and grounded U shape housing is limited by the capacitance between high voltage needle and high voltage electrode and shield. Consequently the quantity of generated ions is limited. For the materials with a high or a very high static charge this ionizer is not efficient, because all the ions generated in corona discharge are attracted by the charge of material, but the charge is still not entirely balanced this way.\nIn Romanian patent no.96753, 1988, inventor Dan D.C. Botez, is presented an apparatus meant for static charge neutralization consisting of a high voltage AC power supply which energizes by a high voltage cable/s one ionizer or more. Referring strictly to ionizer, it is made in the form of a metallic shape supporting insulating modules. A high voltage needle is introduced through each insulating module and makes contact to the conductor of a high voltage cable, placed in the inner part of the insulating module. Between the insulating modules, on both sides, are inserted thin steel wires connected by metallic shape to the ground. This way, each high voltage needle is located in the middle of a rectangle, whose corners contain the peaks of the grounded wires. Between the high voltage needles and peaks of grounded wires takes place a corona discharge, that generates positive and negative ions. The materials with a slight electrical charge get discharged by attracting ions of opposite polarity. High static charges get neutralized mainly by small electrical discharges between the peaks of the grounded wires and charged material. The materials having an average electrical charge get neutralized by small electrical discharges to the peaks of grounded wires, as well by attracting ions of opposite polarity. This type of ionizer has a good efficiency for small, average as well for high and very high static charges. The disadvantages of this ionizer are: a slightly increased ozone generation, because of intense corona discharge; in the inner part of the insulating module is a short path between the high voltage needle via joint of insulated modules to the grounded metallic shape, where a small electrical discharge can take place. The dirt built up on the top surface of ionizer can develop, between high voltage needles and grounded wires, a path with a lower surface resistance, and consequently a leakage current can appear. If the ionizer is not regularly cleaned, a small electrical discharge is possible to develop on this path."}
{"text": "1. Field of the Invention:\nThe present invention relates to jewelry in general, and more particularly to movable and interchangeable gem settings.\n2. Background of the Invention:\nIt is most desirable in jewelry setting to attract the eye of public view toward the ornament, usually a gemstone, around which the piece is designed. It is also desirable to select jewelry which matches clothing or which befits an occasion.\nIn order to meet the above objectives, both user and seller of jewelry must maintain a costly, relatively large variety of pieces on hand. The cost is further increased if matching pieces are desired, say ring and pendant or bracelet, and the gemstone is of the precious or semiprecious type.\nThe present invention is provided to overcome those obstacles."}
{"text": "Object-oriented programming (OOP) is a programming paradigm that uses objects to design software programs. An object may comprise a discrete set of data structures with a defined set of data fields, methods (functions or procedures), properties and events. Each object is capable of receiving messages, processing data, and sending messages to other objects. In other words, each object can be viewed as an independent “machine” with a distinct role or responsibility. As such, a software program may represent a collection of interacting objects.\nObjects for a software program may be executed in a same process or different processes. A process is an instance of a software program that is being executed. Depending on an operating system, a process may be implemented as one or more threads of execution that execute instructions from the software program concurrently. In some cases, objects for a software program may be executed by different threads of a single process, or by different processes, on a single device. In other cases, objects for a software program may be executed by different processes on different devices, such as in a distributed computing environment.\nObjects executed in different processes, or in some cases different threads of a single process, typically utilize some form of inter-process communication (IPC) to exchange data between the different processes. IPC techniques may include operations related to message passing, synchronization, shared memory, and remote procedure calls (RPC). Different IPC techniques may be used based on such factors as bandwidth and latency of communication between processes or threads of a process, and a type of data being communicated.\nA class of IPC techniques may include RPC techniques. RPC techniques may be used to implement a programming concept referred to as “remoting.” Remoting techniques may be used when different objects are implemented by different processes on different devices. Remoting allows a local object executing by a process on one device to access a remote object executing by another process on a different device. For instance, OOP and remoting techniques may be used as a programming model for web applications, such as Rich Internet Applications (RIAs). RIAs are a rapidly growing segment of the software industry. RIAs are web applications that have many of the characteristics of desktop applications, typically delivered either by way of a site-specific browser, via a browser plug-in, or independently via sandboxes or virtual machines. Users generally need to install a software framework using a computer operating system before launching an RIA. The software framework typically retrieves, updates, verifies and executes a given RIA.\nWeb applications such as RIAs utilize various types of remoting techniques. For instance, an RIA may implement objects in a browser as part of a user interface for a web service or web application executing objects on a server. A local object executed by the browser may remotely access certain features of a remote object executed by the web service or web application, such as passing data to the remote object for processing by a method of the remote object. However, conventional remoting techniques are typically limited as to which features of remote objects are accessible, and further, how to access features of remote objects in an efficient manner. It is with respect to these and other considerations that the present improvements are needed."}
{"text": "The production of high pressure laminates is well known, having been carried out for many years. Generally, the procedures for forming such laminates involve providing phenolic (i.e., phenol-formaldehyde) resin impregnated paper core sheets and melamine (i.e. melamine-formaldehyde) resin decorative and overlay sheets, and pressing the stacked resin impregnated sheets under heat on the order of 230.degree.-310.degree. F and under pressure of approximately 800-1600 psi until the resins have become thermoset, thereby providing an extremely hard, attractive and permanent surfacing material known as a \"high pressure laminate\" meetimg NEMA Standards. These high pressure laminates have, for many years, found use as counter and table tops, bathroom and kitchen work surfaces, furniture and cabinet surfacing, wall paneling and partitionings, doors, etc.\nMost general purpose decorative high pressure laminates, of about 1/16 inch thickness, are formed of an assembly comprising a top ply of .alpha.-cellulose paper, about 20-30 lbs. ream weight, impregnated with a partially cured water solution of melamine-formaldehyde condensate; a print ply therebeneath, normally also an .alpha.-cellulose paper, pigment filled, with or without decorative printing on the surface thereof, ranging in weight from 50 to 125 lbs. ream weight and also impregnated with a partially cured amino resin condensate, usually the same melamine resin as used in the overlay ply, and therebeneath a plurality (e.g. six) core plys which are normally 100-130 lbs. ream weight Kraft paper, impregnated with a water or alcohol soluble partially cured phenolic resin condensate. All of such partially cured resin condensates are referred to as being in the \"B-stage\"; in this stage they are thermoplastic and will flow under heat and pressure during the high heat, high pressure laminating procedure.\nThe above described assembly is repeated, back-to-back and front-to-front using separators therebetween, each two laminates being separated by a pressing plate die, until a \"book\" containing, typically, 10 assemblies is obtained. The book is inserted in a high pressure press, and the laminates are cross-linked under the conditions of heat and pressure indicated above. During the pressing cycle, the resins flow and then cure, consolidating the individual plys and forming a comprehensive and infusible cross-linked product. After cure, the book is removed from the press and the laminates are separated, trimmed, and their backs or bottom surfaces sanded to improve adhesion for subsequently gluing to various substrates.\nIn the early years, generally only glossy surface laminates were produced (i.e. using polished molding surfaces), though some attempts were made to provide laminates with some slight surface irregularities to give satin or \"textured\" finishes. In more recent years it has become increasingly desirable to provide various types of irregular surfaces, ranging from relatively shallow depressions such as textured (e.g., satin, matte or semi-gloss) surfaces, to relatively deeply \"sculptured\" or embossed surfaces which have a noticeable three-dimensional effect, such as wood grain, leather, slate, abstract patterns, creative designs, etc. Particularly with the rising costs of natural products, it has become more desirable to provide deeply sculptured high pressure laminates which simulate natural materials such as wood, leather and slate.\nToday, high pressure decorative laminates meeting NEMA Standards are manufactured with smooth glossy surfaces, with textured surfaces, or with deeply sculptured or embossed surfaces. In the latter case, a textured surface is usually superimposed on the embossed surface, and these textured surfaces are normally obtained by utilization of a release, parting or separating membrane or sheet of known type which is located, during the laminating operation, between the upper surface of the uppermost laminate and the molding plate die. In the production of deeply sculptured surfaces which simulate natural products such as slate, leather and wood, it is particularly desirable that the textured surfaces be superimposed over the face of the simulated product, since this provides a softer and more naturally appearing product.\nIn the manufacture of deeply embossed or sculptured surfaces having surface depth variations which duplicate natural products such as slate, leather, wood and the like, it has in the past generally been necessary to use either extremely expensive etched steel or stainless laminating die plates or, alternatively, thermoset pressing plate dies such as described in the Michaelson et al U.S. Pat. Nos. 3,303,081 and 3,311,520. While the use of such dies provides the necessary embossing,, many simulated natural products, such as wood grain, must have appropriate color in registry with the surface embossing or debossing. Unfortunately, the problem of registration of color and embossing has proven very difficult, conventional registry techniques such as used in the printing industry being unreliable on a consistent basis in the manufacture of high pressure laminates, and also being extremely expensive, often involving complex electronic sensing equipment.\nOther methods of providing laminates having registered color and texture have utilized valley printing, namely the inking of high portions of an engraved plate and pressing into the resin while curing. This operation, as well as the mechanical registration of an embossing roll or plate with a print, is difficult to carry into production and/or does not always give a uniform product. In addition, inks may not be compatible with or as durable as the base resin in which case the product will not meet NEMA Standards.\nIn more recent years, a number of other techniques have been developed in an attempt to solve the problems of registration of color and embossing in the manufacture of deeply sculptured high pressure laminates. Certain of these techniques have been found to be highly useful in the preparation of certain specific types of configuration, e.g. see U.S. Pat. Nos. 3,700,537 and 3,698,978. Other recent techniques have been more versitile from the viewpoint of the design of the product (note, e.g., U.S. Pat. Nos. 3,814,647; 3,732,137; 3,802,947 and 3,661,672), but some of these recent techniques have been less than fully desirable because they require special materials and/or are difficult to carry out successfully with the result that the products are sometimes inconsistent or the manufacturing operation is more expensive than is desirable.\nIn copending application Ser. No. 528,776*/(incorporated here by reference) there is disclosed a method of providing registered color and embossing in which the overlay sheet is impregnated with melamine resin containing dispersed pigment, and this overlay sheet is then used in the conventional lay-up over a conventional melamine impregnated print sheet of contrasting color. While the process exemplified in application Ser. No. 528,776 produces a laminate of exceptional appearance particularly when the product is a mottled pattern in which the darker color overlies the lighter color, e.g. simulated leather, such exemplified process does have certain disadvantages. First, it is difficult to carry out and therefore requires exceptional care to impregnate the overlay sheet, without having longitudinal streaks of pigment concentrations. This is due to the mechanics of impregnation and is not a fundamental problem.\nSecondly, and more fundamental is the fact that the process exemplified in Ser. No. 528,776 provides a product which may not have as great an abrasion resistance as is desired. Because much of the pigment is concentrated near the the overlay surface, and the finished laminate surface is heavily embossed, the peaks show loss of color earlier than desirable in the simulated wear tests unless this wear is compensated by usng extraordinary measures, e.g. an extra impregnated overlay sheet, i.e. two overlay sheets. While the use of the second overlay sheet considerably improves the abrasion resistance, this solution reduces the crispness of the color contrast in the fine embossing details with the result that the product is sometimes not as attractive as is desired.\nA third disadvantage of the technique exemplified in Ser. No. 528,766 is that by incorporating the dispersed pigment in the overlay by means of impregnation including first saturating the overlay with pigment dispersed resin and then wiping excess from its surface, the saturating characteristics of the overlay paper have undue influence over the color uniformity of the final product. That is, the normal density and wicking properties of the paper that vary from point-to-point in a sheet may cause variation in resin-pigment concentration at different locations across the area of the sheet. In a conventional overlay, this does not cause problems because there is enough resin to provide suitable physical properties even in low concentration areas, and the variation is not visible because the resin is almost transparent. When the pigmented resin is used, however, concentration variations sometimes result in a mottled appearance in the finished laminate if the pigment is not strongly opaque and, therefore, the technique exemplified in Ser. No. 528,776, while giving a highly realistic mottled simulated product such as leather, is not as suitable for other designs containing larger smooth areas where color variation is not desirable, if the pigment is not strongly opaque."}
{"text": "This invention relates to a new and improved rotary pump for use in many diverse pumping applications, especially in any situation where it would be advantageous to be able to replace or dispose of certain components of the pumping mechanism that have expended their useful lives without having to also dispose of other parts that have additional useful life and/or are more expensive to replace. These concerns are especially relevant in applications employing magnets, which are relatively expensive, as the means of driving the rotation of the pump.\nThe art is replete with examples of rotary pump mechanisms employing various means for driving the rotation of the pump, including but not limited to magnets, gears, electrical conduction and induction, belts and rollers. Pumps employing magnetic driving systems are exemplified by U.S. Pat. Nos. 3,647,324 and 4,037,984 of which Applicant is an inventor. These references disclose pumps for use with biological fluids, especially blood. These pumps are especially useful in heart-lung machines and devices for assisting ailing heart muscles. The pumps are comprised of a housing in which an impeller and a magnet are incorporated. The magnet is magnetically connected to a second magnet outside of the housing. The second magnet is in turn connected to a motor. When the motor rotates the second magnet, the first magnet also rotates, turning the impeller with it.\nOther pumps incorporating similar magnetic driving mechanisms are represented by U.S. Pat. Nos. 3,306,221, 3,321,081, 3,575,536, 4,135,253 and 4,165,206.\nThese magnetic systems keep the pumped fluid separated from the mechanical driving elements of the pump, for example the motor. This feature provides the advantages of (1) preventing fluids from contacting the driving mechanism, thus preventing corrosion, electrical shorts, seizing and other undesirable results that render the driving mechanism useless; (2) preventing lubricants, detritus and other foreign material from contacting the fluid being pumped, especially when the fluid is blood; and (3) preventing the fluid being pumped from contacting air, especially when the fluid is blood.\nTo maintain separation between the fluid and the driving mechanism, many devices in the prior art incorporate an expensive driving mechanism, for example magnet-driven shafts or other rotatably components, into the impeller section of the pumping device. Since components associated with rotation are usually the mechanical elements that wear out most quickly, replacement of these components means replacement of the entire impeller section including components in that section which may not have outlived their useful lives. Alternatively, in some applications, especially those in the medical field where the pump is used in contact with bodily fluid or preparations of body tissues, it may be necessary to dispose of the impeller section of a pump because it cannot be reused due to the possibility of contamination or other undesirable reaction between substances from one patient and those of another. With pumps of the prior art, if the impeller and other elements contacting the biological fluid must be discarded, other components, including parts of the magnetic or other driving mechanism, must be unnecessarily discarded with the impeller. The unnecessary, inefficient and uneconomical loss of components that would still be useful is a problem that is common to many rotary pumping devices in the prior art.\nIt would be desirable, therefore, to provide a rotary pump that keeps fluid away from the driving mechanism, yet allows replacement or disposal of components of the pump which cannot be reused or continue to be used, without also disposing of components that are still useful."}
{"text": "Temperature and humidity sensors capable of detecting temperatures and humidity are disclosed in, for example, Japanese Unexamined Utility Model Registration Application Publication No. 7-20536 (Patent Document 1), Japanese Unexamined Patent Application Publication No. 6-148122 (Patent Document 2), and Japanese Unexamined Patent Application Publication No. 58-70153 (Patent Document 3). These temperature and humidity sensors include a moisture sensitive film for humidity measurement and detect humidity based on a change in magnitude of resistance or electrostatic capacity.\nFurther, detection of electrostatic capacity of a humidity sensor using an oscillation circuit is disclosed in Japanese Unexamined Patent Application Publication No. 61-14556 (Patent Document 4) and Japanese Unexamined Patent Application Publication (Translation of PCT Application) No. 2012-508877 (Patent Document 5), for example.\nPatent Document 1: Japanese Unexamined Utility Model Registration Application Publication No. 7-20536.\nPatent Document 2: Japanese Unexamined Patent Application Publication No. 6-148122.\nPatent Document 3: Japanese Unexamined Patent Application Publication No. 58-70153.\nPatent Document 4: Japanese Unexamined Patent Application Publication No. 61-14556.\nPatent Document 5: Japanese Unexamined Patent Application Publication (Translation of PCT Application) No. 2012-508877.\nIt can be considered to configure an oscillation circuit as disclosed in Patent Documents 4 and 5 using capacitance and resistance of a humidity sensor based on an electrostatic capacity technique as disclosed in Patent Documents 1 through 3 as well as inductance of an inductor that is provided separate from the humidity sensor.\nIn the case where an oscillation circuit is a CR oscillation circuit, because, in general, the phase rotation is gentle and oscillation frequency precision is low, and the circuit is susceptible to parasitic capacitance and resistance components, electrostatic capacity of a humidity sensor is required to be large. Accordingly, in a humidity sensor in which a moisture sensitive film is provided between two electrodes like the sensors in Patent Documents 1 and 2, areas of the electrodes need to be enlarged, thickness of the moisture sensitive film needs to be thinned, and so on. However, miniaturization of the sensor becomes difficult if the areas of the electrodes are enlarged, and an amount of change in electrostatic capacitance becomes small and the manufacturing of the sensor becomes difficult if the thickness of the moisture sensitive film is thinned. In a humidity sensor provided with linear electrodes like the sensor in Patent Document 3, intervals between the electrodes need to be narrowed. However, narrowing the intervals makes it difficult to manufacture the sensor."}
{"text": "The background description includes information that may be useful in understanding the present invention. It is not an admission that any of the information provided herein is prior art or relevant to the presently claimed invention, or that any publication specifically or implicitly referenced is prior art.\nMotors with a rotor having permanently magnetized poles and a stator having electrically energized field poles suffer numerous disadvantages. The current for such motors is usually supplied from an alternating current source or from communicators rotating with the rotor. However, the maximum speed of such motors is limited by the frequency of the alternating current or by the ability to rapidly reverse the flow of the current in the field coil.\nTo address these problems, U.S. Pat. No. 2,968,755 to Baermann discloses a motor, which includes stator poles with permanent magnet means for magnetizing each stator pole. Baermann's stator poles also have a remotely actuated magnetic means of a greater magnetic strength than the permanent magnetic means, which could be used to reverse the magnetic polarity of its corresponding stator pole without needing to provide reversing alternating current. However, Baermann's stator only makes use of magnetic flux from one pole of each permanent magnet, and reversing the magnetic polarity of Baermann's stator is energetically demanding.\nIn U.S. Pat. No. 5,825,113, Lipo et al. discloses permanent magnetic machines that employ a field winding that can be used to weaken or boost an existing magnetic field. Lipo's motors comprise a pair of arched permanent magnets embedded in a stator yoke, a field winding, and armature windings. Half of Lipo's stator poles are dedicated magnetic north poles and half are dedicated magnetic south poles diametrically opposing the dedicated magnetic north poles. U.S. Pat. No. 2,816,240 to Zimmerman also employs similar principles to those disclosed by Lido. However, both Lipo and Zimmerman only use field windings to weaken or boost the variable magnetic fields.\nAll publications identified herein are incorporated by reference to the same extent as if each individual publication or patent application were specifically and individually indicated to be incorporated by reference. Where a definition or use of a term in an incorporated reference is inconsistent or contrary to the definition of that term provided herein, the definition of that term provided herein applies and the definition of that term in the reference does not apply.\nThus, there is still a need for permanent magnet offset systems that vary the strength and quality of magnetic fields in a more energy-efficient manner."}
{"text": "1. Field of the Invention\nThe present invention relates to an electronic track lighting system, and in particular, to a lighting system which lights a lamp load such as a halogen lamp at a high frequency by an inverter circuit using an electric power fed via a power track.\n2. Description of the Related Art\nConventionally, an inverter circuit disclosed in Japanese Patent No. 2770337 has been well known as an inverter circuit for lighting a halogen lamp, and the disclosed circuit has a configuration as shown in FIG. 17. In the inverter circuit, a portion surrounded by a broken line shown in FIG. 12 is provided as a unit. The inverter circuit is built in a transformer unit A and is used together with a lamp load 3 in the form as shown in FIG. 13. As shown in FIG. 14, the electronic transformer unit A is connected to a power track 6 mounted on a ceiling via a power track connective portion 4 provided on an upper portion thereof. A commercial alternating current voltage of AC 100V is supplied to the electronic transformer unit A through the power track 6. A lamp unit 5 suspended from a lower portion of the electronic transformer unit A is mounted with a lamp load 3 such as a halogen lamp or the like.\nFor example, such a luminaire (lighting equipment) is used for showcase (display case) lighting, light up of show windows, and it is strongly desired that the electronic transformer unit A connected to the power track 6 is miniaturized. However, the prior art employing the aforesaid 100V power track system has the following problem. More specifically, as seen from a circuit shown in FIG. 15, each of a plurality of the electronic transformer units A includes a filter circuit for reducing a power supply feedback noise and a protective circuit for protecting a circuit when a lightning surge is applied, respectively. Therefore, size of the electronic transformer unit A becomes large.\nA 12V power track system as shown in FIG. 16 is often used in general, and is constructed in the following manner that a secondary side (output side) of an output transformer T1 of the electronic transformer unit A is connected to the power track 6, and a plurality of 12V mini-halogen lamps are connected to the power track 6 via a lamp unit to be lit. In this case, the power track 6 is supplied with a high frequency voltage of 12V transformed from a commercial supply voltage of AC 100V by the electronic transformer unit A.\nMoreover, in U.S. Pat. No. 5,180,952, proposed is a method of supplying, to the power track system, a voltage which is converted from a commercial power supply voltage into a high frequency voltage to downsize each unit.\nIn the AC 100V power track system as shown in FIGS. 14 and 15, the electronic transformer unit A must be connected one by one with the power track. Therefore, depending upon a size of the electronic transformer unit A connected to the power track 6, there is a limitation of the number of electronic transformer units A capable of being connected to a predetermined length of the power track. Conversely, comparing to the 12V power track system, a current flowing through the power track 6 is less, and thus the AC 100V power track system has a merit such that a voltage drop by an impedance of the power track 6 is smaller than the 12V power track system and does not affect brightness of the light unit.\nOn the other hand, in the high frequency 12V power track system shown in FIG. 16, one electronic transformer unit A is shared with a plurality of halogen lamps 3. Therefore, each halogen lamp 3 is connected to the power track 6 by only lamp unit 5. As a result, the high frequency 12V power track system has merit such that the whole of lighting system can be miniaturized. However, as described above, a voltage is low, that is, 12V, and thus a current flowing through the power track 6 is 9 times as much as the AC 100V power track system. Therefore a voltage drop by an impedance component of the power track 6 is large, and when the long power track is used, the lamp unit far from the transform unit A becomes dark. Moreover, in attaching/detaching the lamp unit 5 to/from the power track 6, there is the following problem. Since a current flowing through a connective fitting 4 with the power track 6 is large, an arc is generated at the attaching or detaching.\nThe present invention has been made in view of the above problems in the prior art. It is therefore an object of the present invention to provide a lighting system which can suppress a current flowing through a power track so that no arc is generated between a connective metal fitting and the power track when attaching and detaching a lamp unit."}
{"text": "Conventional metal cans which are provided with a bail to carry same, such as for example one gallon paint cans, are provided with ears which are welded, clinched or otherwise attached, on diametrically opposed sides of the side wall of the can whereby to attach the bail. In the fabrication of the can the application of the ear to the paint can causes problems in the assembly process. For example, when metal ears are welded to the side wall of the can the welding process can cause some destruction and cause bare metal to be exposed on the inside of the can. When latex paint is placed in such cans, consumers can find rust floating on top of their paint due to its contact with bare metal. The precise ear placement is also problematic with welded ears, often leading to rejecting products and requiring extra inspection.\nA further problem with paint cans on which ears are welded is that when labels are applied to the can the machine must take into consideration the location of the ears and this poses a restriction on the speed of operation. It also causes defective label applications and again resulting in rejection of cans where labels are improperly applied and therefore requiring more labor to recycle such cans. The current method of applying labels to the paint can is literally a slow rolling process. Before the label is applied, the can has to be indexed so that the ears are properly timed as they pass through the labeling process, a definite bottleneck. Once indexed, the can is rolled over glue and then over a paper label."}
{"text": "When two devices are incompatible or lack access to a connecting network, data transfer between them might still be possible using displayed and captured visual codes such as a QR code, a PDF417 code, etc.\nHowever, the standards for visual codes employ maximum data limits which are smaller than the size of many image files or files containing non-image biometric data, which are often used to prove identity, e.g., to law enforcement officials or airlines."}
{"text": "On-line transactions between consumers and merchants on the World Wide Web (WWW) are becoming increasingly more numerous as the public becomes more facile in making purchases on the Internet. Such transactions can be for the purchase of \"soft\" goods, i.e., information, software and other material available in electronic form that can be delivered in real time to a user's client terminal. Such transactions can also be for the purchase of conventional \"hard\" goods, where the purchased merchandise are delivered off-line. Conventional on-line payment options generally involve the use of credit cards wherein the user provides his or her credit card number on-line or off-line to the merchant provider to pay for the \"hard\" or \"soft\" purchase to be delivered on-line or off-line.\nWhere the transactions involve a relatively small cost, for example $10 or less, the credit card system of payment is too expensive. Further, the credit card payment system excludes potential customers who do not have a credit card, or those who do but do not \"trust\" either providing their credit card number on-line, or do not want to use their credit card for such purchases. It would be advantageous, therefore, that some trusted transaction intermediary perform the functions of authenticating a user on the WWW and authorizing the transaction. Once such a transaction intermediary authenticates the user and authorizes the transaction, the merchant is alerted to provide the goods or services which are the subject of the transaction and an account associated with the user is billed for the transaction amount. Advantageously, once the user has registered with the transaction intermediary, no sensitive billing information, such as a credit card number, needs to be sent to the merchant.\nIn co-pending application Ser. No. 08/532,336 filed Sep. 22, 1995 by Y. Ronen, co-inventor herein, and assigned to the same assignee as the present application, a method of billing for services and/or goods ordered over the Internet from a merchant is disclosed. As described therein, the customer/user places a real or a virtual telephone call to the merchant's 900 telephone number and is charged by the telephone company an amount for that call that is representative of the cost of the goods or services being ordered via the Internet from the merchant's server. The merchant's server associates the 900-number telephone call with the request via the Internet for the goods or services in order to authenticate and authorize the transaction. In co-pending patent application Ser. No. 08/636,109 filed Apr. 22, 1996, co-invented by the same Mr. Ronen who is co-inventor herein, and assigned to the same assignee as the present application, a method of billing for transactions conducted over the Internet is disclosed in which a billing server connected on the Internet receives the IP address assigned to that user's client for the current user session and an indication of the user's identity from the user/customer's Internet Access Provider (IAP). In response to a chargeable transaction, the merchant's server transmits to a billing platform the IP address identity of the user making the transaction and the cost associated with the transaction. The billing server then cross-references the IP address associated with the cost of the transaction received from the merchant's server with the IP-address/user-identity relationship received from the IAP to properly charge an established account of the user for the transaction. Such an account is established by the user prior to the execution of the transaction for billing in a predetermined manner to, for example, the user's selected credit card, the user's debit card, the user's telephone account associated with his or her phone number, the user's merchant credit card, or other billing mechanism. For this billing methodology, arrangements thus need be established between billing server and the large number of different Internet Access Providers that provide Internet access to a tremendously large customer base since for each user's session the IAP must be programmed to forward to the appropriate billing server the relationship between the user's currently assigned IP address and identity.\nA billing methodology that minimizes the steps that need be performed to obtain authorization and approval for an Internet transaction is therefore desirable. Further, a billing methodology that requires interactions between only the user, the merchant, and the provider of the billing service is advantageous."}
{"text": "The following generally relates to a method for controlling an electrical tool.\nElectrical tools, such as an electrical torsion spanner, electrical trimmer, and electrical pruner, perform specific operations in a way that a motor drives a function member to move in a reciprocating manner. As far as this type of electrical tool is concerned, the tool is expected to reset automatically when the tool function is not accomplished for the sake of security purposes and the user's convenient operation. Take the electrical torsion spanner 100 shown in FIG. 1 as an example, the electrical torsion spanner 100 is an open-ended spanner. When using this tool, the user hopes that the spanner is placed in a state shown in FIG. 1 upon initial use each time. However, after power off, the electrical torsion spanner 100 is usually in a random position as shown in FIG. 2 and therefore causes great inconvenience to the user's use thereof."}
{"text": "1. Field\nExample embodiments relate to an automatic packing apparatus which automatically packs products in boxes and an automatic packing method using the same.\n2. Description\nIn general, after various products are manufactured, the products are packed in boxes in order to prevent the products from being damaged and to ship/handle the products.\nThe packing boxes are made of a material, such as thick corrugated cardboard, so as to protect the products from external impact. The manufactured packing boxes are stored in an unassembled state in which the boxes are unfolded in a flat shape so as to reduce a space occupied thereby before being used to pack the products. Then, when the packing boxes are used to pack the products, the packing boxes are assembled (folded) to pack the products.\nA worker folds a packing box into an original shape, and then covers a product with the packing box, thereby packing the product in the packing box. Here, the product is covered with the box downwardly from above. Packing of a large-sized home appliance product, such as a refrigerator, a TV or an air conditioner, in such a box requires at least two workers.\nAs described above, workers manually pack products in boxes and thus packing requires considerable manpower, and a large amount of time is consumed packing each product and thus productivity is reduced."}
{"text": "1. Field of the Invention\nThe present invention relates to a plasma processing apparatus, and more specifically, relates to a plasma processing apparatus capable of suppressing the occurrence of particles caused by reaction products.\n2. Description of the Related Art\nIn the field of manufacturing semiconductor devices, plasma processing apparatuses using plasma are adopted as means for realizing microfabrication to respond to the increasing demands for highly precise high-speed microfabrications (such as dry etching), to realize enhanced performance of semiconductor integrated circuit elements and cut down costs. Among such arts, the art of dry etching is often applied, in which processing gas is used to process films deposited on the sample and evacuate the residual gas in the processing chamber via a turbo-molecular pump or the like to the exterior. Various kinds of gases (such as Ar, Cl2 and HBr) can be used as the processing gases, but some gases are not sufficiently evaporated, depending on the kind of films deposited on the sample being processed. In such cases, gases are not sufficiently discharged and become the cause of reaction products, which deposit on the inner wall of the vacuum processing chamber and on the inner wall of an inductive vacuum processing chamber lid disposed on the upper portion of the vacuum processing chamber. When excessive reaction products are deposited on the inner wall of the vacuum processing chamber and the inner wall of the vacuum processing chamber lid, deposition films may fall off from the surface of the inner wall and attach to the surface of the sample as particles.\nOne type of prior art plasma processing apparatuses is an inductive plasma processing apparatus having a high frequency power supply connected to a coil-like induction antenna disposed on an outer circumference of a vacuum processing chamber lid for supplying high frequency power to the coil-like induction antenna, to thereby generate plasma. The drawback according to this type of plasma processing apparatus adopting an induction antenna is that the bonding state of the induction antenna and the plasma within the vacuum processing chamber lid is varied by the reaction products deposited on the inner wall of the vacuum processing chamber and the vacuum processing chamber lid, by which the etching rate, the uniformity of etching, the perpendicularity of the etching profile and the state of deposition of reaction products on the etched side walls is varied with time.\nNext, Japanese patent application laid-open publication No. 2004-235545 (patent document 1) discloses the following method as the method for preventing reaction products from depositing on the inner wall of the vacuum processing chamber lid. A Faraday shield capacitively coupled with plasma is arranged between induction antennas disposed at the outer circumference of a vacuum processing chamber lid and plasma, and high frequency power is supplied to the Faraday shield to generate self bias to the inner wall of the vacuum processing chamber lid, drawing ions in the plasma by the electric field within the ion sheath and sputtering the surface of the inner wall of the vacuum processing chamber lid. The sputtering is used to suppress or remove the deposition of reaction products on the inner wall of the vacuum processing chamber lid, and perform cleaning of the inner wall of the vacuum processing chamber lid. In the present specification, the vacuum processing chamber lid is airtightly fixed to the upper portion of the vacuum processing chamber.\nAs described, a self bias is generated on the surface of the inner wall of the vacuum processing chamber lid via the high frequency voltage applied to the Faraday shield, by which the electric field of the ion sheath is increased and the energy of ions reaching the inner wall of the vacuum processing chamber lid is enhanced. Therefore, ion sputtering via ions becomes prominent, enabling more reaction products to be cleaned. Further, most appropriate cleaning becomes possible by controlling the high frequency voltage supplied to the Faraday shield.\nFurther, the following method is disclosed in Japanese patent application laid-open publication No. 2005-259836 (patent document 2) as a method for cleaning the whole surface of the inner wall of the vacuum processing chamber lid in a uniform manner. Patent document 2 discloses a plasma etching apparatus having a Faraday shield divided into multiple parts and providing independent high frequency power supplies respectively to each of the divided Faraday shields, to thereby enable control of the voltages applied to the respective Faraday shields.\nAccording to the disclosed plasma processing apparatus, in order to remove the reaction products generated from samples during etching and thickly deposited on a center portion of the vacuum processing chamber lid (top plate) closest to the sample surface, a high frequency power which is higher than the power supplied to the outer circumference portion of the top plate is supplied to the divided Faraday shield disposed at the center portion of the top plate. According to this arrangement, reaction products deposited both on the center portion and the outer circumference portion of the top plate can be cleaned appropriately.\nFIG. 10 shows the state of deposition of reaction products attached to the inner wall of the vacuum processing chamber lid in a plasma processing apparatus. When etching is performed, a large amount of reaction products is generated from the sample being processed, and the generated reaction products are attached to and deposited on the inner wall of the vacuum processing chamber lid (hereinafter referred to as a window 12) and the inner wall surface of a vacuum processing chamber 2 as a film (18).\nIn contrast, at an outermost circumference portion (portion B) of the inner wall surface of the window 12, the self bias is reduced since the vacuum processing chamber 2 is grounded, so that the energy of ions with respect to portion B is reduced. Ion sputtering rarely occurs in portion B, and therefore, reaction products remain in portion B. As a result, the inner wall surface of the window 12 is divided into a region (portion A) where the reaction products are completely removed, and a region (portion B) where reaction products easily remain.\nIn order to eliminate residual reaction products in portion B, it is necessary to set the high frequency voltage applied to the Faraday shield 17 (Faraday shield voltage: hereinafter abbreviated as FSV) as high as possible to remove the reaction products deposited on the inner wall of the window 12. According to this method, however, the inner wall surface of the window 12 of portion A is chipped excessively.\nAs described, if the inner surface of the window 12 is chipped and the speed of consumption of the window 12 itself is increased, the replacement cycle of the window 12 is shortened. Further, the chipped window 12 becomes the source of reaction products and re-deposits on the wafer surface, causing instability of etching rate of the process performance and in-plane unevenness of the wafer etching rate and the like. Another drawback is that the chipped surface of the window 12 turns into particles and attach to the sample surface. Therefore, in order to overcome the above-mentioned problems, it is necessary to remove the reaction products deposited on the whole inner wall surface of the window 12 in a most suitable manner.\nPrior arts such as that disclosed in patent document 2 enables to remove the reaction products deposited on the whole inner wall surface of the window 12 in a suitable manner, but since independent high frequency power supplies must be disposed in response to the divided Faraday shields and antennas, the application thereof to industrially-applied plasma processing apparatuses causes increase of costs and requires a large mounting space.\nFurthermore, there is a drawback that cannot be solved by the prior art disclosed in patent document 2. This drawback will be described with reference to FIG. 5. Since the present drawback is not related to the division number of Faraday shields, the plasma processing apparatus illustrated in FIG. 5 has a Faraday shield that is not divided into multiple parts.\nIn FIG. 5, the dotted lines with arrows show high frequency currents flowing from a Faraday shield 17 into the vacuum processing chamber 2 during plasma processing. The Faraday shield 17 is capacitively coupled with plasma, and the high frequency current from the Faraday shield 17 flows through the plasma as shown by the arrows into the vacuum processing chamber 2 being grounded.\nNow, when the high frequency voltage applied to the Faraday shield 17 is denoted by Va, the area of the Faraday shield 17 is denoted by Sa, the high frequency voltage generated in the ion sheath on the surface of the inner wall of the vacuum processing chamber 2 is denoted by Vb, and the surface area of the ion sheath is denoted by Sb, the following expression is satisfied, as taught in cited document (Michael A. Lieberman, Allan J. Lichetenberg, ED Research, “Principles of Plasma Discharges and Materials Processing” Chapter 11).\n                                          V            b                                V            a                          =                              (                                          S                a                                            S                b                                      )                    n                                    [                  Expression          ⁢                                          ⁢          1                ]            \nExpression 1 shows that a self bias is not only generated at the inner wall surface of the window 12 by the high frequency voltage applied to the Faraday shield 17, but also generated at the inner wall surface of the vacuum processing chamber 2 into which high frequency current flows.\nTherefore, the inner wall of the vacuum processing chamber 2 is sputtered via the ions accelerated by Vb. As described, when the inner wall of the vacuum processing chamber 2 is sputtered, the lifetime of the respective parts composing the interior of the vacuum processing chamber 2 is reduced, and the frequency of replacement of the parts is increased. Thus, the running costs of the plasma processing apparatus cased by replacement of parts is increased. Even further, the chipped inner wall of the vacuum processing chamber 2 becomes the source of particles, by which the mass production performance of the plasma processing apparatus is deteriorated."}
{"text": "The invention relates to compounds and compositions useful for treating disorders related to KIT and PDGFR.\nThe enzyme KIT (also called CD117) is a receptor tyrosine kinase expressed on a wide variety of cell types. The KIT molecule contains a long extracellular domain, a transmembrane segment, and an intracellular portion. The ligand for KIT is stem cell factor (SCF), whose binding to the extracellular domain of KIT induces receptor dimerization and activation of downstream signaling pathways. KIT mutations generally occur in the DNA encoding the juxtumembrane domain (exon 11). They also occur, with less frequency, in exons 7, 8, 9, 13, 14, 17, and 18. Mutations make KIT function independent of activation by SCF, leading to a high cell division rate and possibly genomic instability. Mutant KIT has been implicated in the pathogenesis of several disorders and conditions including systemic mastocytosis, GIST (gastrointestinal stromal tumors), AML (acute myeloid leukemia), melanoma, and seminoma. As such, there is a need for therapeutic agents that inhibit KIT, and especially agents that inhibit mutant KIT.\nPlatelet-derived growth factor receptors (PDGF-R) are cell surface tyrosine kinase receptors for members of the platelet-derived growth factor (PDGF) family. PDGF subunits-A and -B are important factors regulating cell proliferation, cellular differentiation, cell growth, development and many diseases including cancer. A PDGFRA D842V mutation has been found in a distinct subset of GIST, typically from the stomach. The D842V mutation is known to be associated with tyrosine kinase inhibitor resistance. As such, there is a need for agents that target this mutation."}
{"text": "The incidence of gram negative bacteria in the United States has been estimated to be approximately 100,000 to 300,000 cases per year, with a mortality rate of 30-60%. Antibiotics are commonly used as the primary chemotherapy for this disease; however, their bactericidal action can result in disruption of the bacterium and concomitant release of endotoxin, i.e., the lipopolysaccharide (LPS) moiety of the bacterial outer membrane. The liberated LPS induces a number of pathophysiological events in mammals (collectively referred to as gram-negative endotoxemia or sepsis syndrome). These include fever, generalized inflammation, disseminated intravascular coagulation (DIC), hypotension, acute renal failure, acute respiratory distress syndrome (ARDS), hepatocellular destruction and cardiac failure.\nIt has previously been established that, for infections caused by gram-negative bacteria, sepsis is related to the toxic components of the bacteria. Specifically, among the well-described bacterial toxins are the endotoxins or lipopolysaccharides (LPS), a cell-wall structure of the gram-negative bacteria. These molecules are glycolipids that are ubiquitous in the outer membrane of all gram-negative bacteria. While the chemical structure of most of the LPS molecule is complex and diverse, a common feature is the lipid A region of LPS (Rietschel, et al., in the Handbook of Endotoxins, 1: 187-214 eds. R. A. Proctor and E. Th. Rietschel, Elsevier, Amsterdam (1984)); recognition of lipid A in biologic systems initiates many, if not all, of the pathophysiologic changes of sepsis. Because lipid A structure is highly conserved among all types of gram-negative organisms, common pathophysiologic changes characterize gram-negative sepsis. It is also generally thought that the distinct cell wall substances of gram-positive bacteria and fungi trigger a similar cascade of events, although the structures involved are not as well studied as gram-negative endotoxin.\nAlthough endotoxin initiates septic shock, it has little or no direct toxic effect on tissues; instead, it triggers an immunobiological response leading to a cascade of release of cytokines such as tumor-necrosis factor (TNF), interleukin-1, interleukin-6 and interleukin-8, and other biological mediators such as nitric oxide, as well as an array of secondary mediators (e.g., prostaglandins, leukotrienes, interferons, platelet-activating factor, endorphins and colony-stimulating factors).\nTherapies for endotoxin-related diseases have generally been directed towards controlling the inflammatory response. Such therapies include corticosteriod treatment, suggested to ameliorate endotoxin-mediated cell membrane injury and to reduce production of certain biological mediators; administration of antibodies designed to neutralize bacterial LPS; treatment with agents to suppress hypotension or with naloxone which apparently blocks the hypotensive effects associated with sepsis syndrome; and treatment with nonsteroidal anti-inflammatory drugs, purported to block cyclooxygenases and thereby decrease the production of certain secondary mediators such as prostaglandins and thromboxane.\nHowever, none of these therapies to date has resulted in significant reduction in the morbidity and mortality resulting from sepsis and septic shock syndrome. Thus there is a long felt need for agents to affirmatively treat this disorder.\nCertain lipodisaccharides are disclosed in Macher et al., Great Britain patent 2,179,945, Meyers et al., Great Britain patent 2,220,211, Shiba et al., European patent 172,581, Anderson et al., U.S. Pat. No. 4,495,346 and Shiba et al., U.S. Pat. No. 5,066,794.\nChrist, et al., “Anti-Endotoxin Compounds,” U.S. Pat. No. 5,530,113, filed Aug. 25, 1992, the contents of which are included by reference, also disclose certain disaccharide compounds, such as B531 shown below, useful for the treatment of endotoxemia.\n\nLipodisaccharides with enhanced pharmacological selectivity, efficacy and increased persistence in action were disclosed in U.S. Pat. No. 5,935,938, the entire contents of which are incorporated herein by reference. Specifically, lipodisaccharide B1287 was identified as having a better pharmacological profile than some of its congeners, and thus shows great promise for the treatment of LPS-mediated disorders.\n\nHowever, existing methods for its preparation typically involve >36 synthetic steps, and are thus not well suited for industrial applicability. Specifically, the synthetic approaches for preparing B1287 that are disclosed in U.S. Pat. No. 5,935,938 are lengthy, hazardous (e.g., involve azide chemistry) and thus inadequate for large-scale syntheses.\nClearly, there remains a need to develop efficient and high yielding synthetic methodologies to access a variety of analogues of Lipid A, particularly B1287 and stereosiomers thereof, which compounds are useful as in the prophylactic and affirmative treatment of endotoxin exposure including sepsis, septicemia, endotoxemia, various forms of septic shock and related disorders using novel liposaccharide analogs."}
{"text": "1. Field of the Technology\nThe technology relates to a developing unit and an image forming apparatus using this. The technology, in particular, directed to a developing unit including an agitating and conveying portion for agitating and conveying a developer in the developer storing portion for storing the developer containing an electrostatically chargeable toner and a magnetic carrier to visualize an electrostatic latent image formed on a photoreceptor drum with the toner, which is constructed to agitate the developer in the developer storing portion without unevenness as well as to an image forming apparatus using this developing unit.\n2. Description of the Prior Art\nIn a conventional image forming apparatus such as a printer etc. using the electrophotography, a printout of an image is formed by electrifying a photoreceptor drum that is driven rotationally by a charger, forming an electrostatic latent image by illuminating the photoreceptor drum with light in accordance with image information, forming a toner image by applying the toner to the electrostatic latent image by a developing unit and transferring the toner image to a sheet material or other recording media.\nAs the developer used for the developing unit of the image forming apparatus, use has been made of, for example a dual-component developer which is prepared by mixing two components, or an electrostatically chargeable toner and a magnetic carrier.\nIn the thus constructed image forming apparatus, the developer is stored in the developer storing portion as a part of the developing unit and is agitated inside the developer storing portion and conveyed to the developing roller by an agitating and conveying means.\nConventionally, the means for agitating and conveying the developer in the above way has been usually realized with conveyor belts, conveying screws, paddles and the like.\nAs an example using a conveyor belt, there has been proposed a configuration for agitating and conveying the toner in the toner supply container, in which an circulating endless conveyor belt having conveying projections on its outer surface and agitating projections on its inner surface is arranged to agitate and convey the toner inside the toner supply container (see Patent document 1: Japanese Patent Application Laid-open Hei 9-138579).\nAlso three has been proposed a configuration for agitating and conveying the developer in the developing unit, in which a magnetic belt is adopted as a conveyor belt so that the magnetic belt carrying the developer on its surface by magnetic force is moved to convey the developer toward the developing roller while agitating the developer therearound (see Patent document 2: Japanese Patent Application Laid-open 2003-302837).\nFurther, there has been proposed another configuration, in which a screw is used to send back the collected toner into the developing unit while inside the developing unit the collected toner is conveyed by a conveyor belt (see Patent document 3: Japanese Patent Application Laid-open 2001-201930).\nHowever, since in any of the above conventional technologies, the developer containing toner is agitated and conveyed by moving the belt surface or agitating and conveying projections of the conveying belt upwards or downwards, the agitating effect in the vertical direction is weak. Further, there is a problem that when the developer is conveyed toward the developing roller, the developer is unlikely to be agitated and dispersed uniformly with respect to the axial direction of the developing roller or across the main scan direction."}
{"text": "It is customary to stow loads on top of the cargo hatches on the weather deck of ships and to cover the loads with a tarpaulin. In hard weather, waves will wash over the ship and the load on the hatches might be damaged."}
{"text": "The present invention is in the field of preparing self-assembled monolayers on a metal using electrolysis.\nMolecular self-assembly of monolayers (SAMs) provides a simple method to form highly ordered two-dimensional organic assemblies. Among the various systems that display this behavior, SAMs formed by the chemisorption of alkanethiols on gold to produce strong gold-thiolate bonds is particularly convenient because of the ease of their preparation. [1-3] Well-ordered SAMs form spontaneously on gold surfaces within a short period upon immersion of the surface in either a dilute solution or the vapor of an alkanethiol of interest. These monolayers have been used as model systems for fundamental studies of wettability [4], adhesion [5], biocompatibility [6], fouling [7], as well as serving as the basis for building multilayers [8] and bio- and analytical sensors [9], for the immobilization of biomolecules [10], and for preparing patterned surfaces [11].\nOne of the most attractive characteristics of self-assembled monolayers (SAMs) is the facility of their preparationxe2x80x94simply exposing a substrate surface to an appropriate adsorbate-[1-3] though this facility prevents selective formation at particular surfaces in the presence of others of the same composition. While contact printing of SAMs has proven useful in preparing mesoscale patterns on various substrates [12], formation of monolayers on only particular features of a pre-existing patternxe2x80x94such as an electrode array, an integrated circuit, or a MEMS devicexe2x80x94has remained an elusive synthetic challenge. Hence, treating an array of identical gold electrodes with a solution of an alkanethiol or alkyl disulfide would result in formation of a monolayer on all of the electrodes. One approach to controlling self-assembly in these systems has focused on the inability of thiolate ions to covalently bind the gold surface directly, without concomitant oxidation. [13] Hence, gold electrodes immersed in a solution of thiolate ions will only adsorb a SAM if held at a sufficiently high potential. Electrochemistry, in this case reduction, can also be used to remove SAMs that had previously been adsorbed. [14] Hence, an alternative strategy for producing patterns involves allowing indiscriminate self-assembly to occur, followed by articulation of a pattern by the selective removal of the SAM from certain substrates.\nAlkylthiosulfates, also known as Bunte salts, can be used to synthesize disulfides by oxidation [15], acidic hydrolysis [16-17], or alkaline degradation [18]. Disulfides also can be formed from Bunte salts electrochemically. [19-20] This method has been extended to form polydisulfides from xe2x80x9cdoublexe2x80x9d Bunte salts, molecules carrying two thiosulfate groups, using electrochemistry with gold electrodes. [21]\nDisclosed is a method of preparing self-assembled monolayers on a metal comprising electrolyzing a thiosulfate compound in a solvent, where the electrolysis is performed at a voltage for a period of time.\nAlso disclosed is a method of preparing self-assembled organic monolayers on a metal comprising (a) contacting said metal with a solution comprising an organic thiosulfate compound, and (b) electrolyzing said organic thiosulfate compound by applying on said metal sufficiently high anodic potential for sufficient time to result in the oxidative self-assembly of said monolayers on said metal.\nFurther, a method is disclosed for the selective formation of self-assembled organic monolayers on a first metal electrode in the vicinity of a second metal electrode, comprising (a) contacting said metal electrodes with a solution comprising an organic thiosulfate compound under conditions such that chemisorption of said organic thiosulfate compound onto said first and second electrodes does not occur, and (b) electolyzing said organic thiosulfate compound by selectively applying on said first metal electrode sufficiently high anodic potential for sufficient time to result in the oxidative self-assembly of said organic monolayers on said first electrode."}
{"text": "In recent years, there is a method to issue a reserved-seat ticket having a seat number printed thereon on the day of the event in order to prevent resale of a ticket of an event and the like. In this method, a user brings a reserved-seat voucher having individual information described thereon (for example, his or her full name and the like) and the user's identification card to the event site. In a case where an administrator of the event determines that the individual information on the reserved-seat voucher and the user's identification card match each other, a reserved-seat ticket is issued to the user. In this method, it is necessary to issue as many reserved-seat vouchers and reserved-seat tickets as the number of users. Therefore, depending on the number of users, the cost of the tickers would be enormous. Such problem is not limited to the tickets, and is a common problem that occurs when multiple types of sheets having information formed thereon are issued to the user."}
{"text": "1. Field of the Invention\nThis invention relates to coal flotation, and more particularly to coal froth flotation utilizing gaseous carbon dioxide for production of clean coal concentrates.\n2. Description of the Prior Art\nFroth flotation is a physicochemical separation process that depends on the attachment of air bubbles to hydrophobic particles. Other (hydrophilic) particles are wetted by the aqueous phase and will not attach to the air bubbles. Thus the separation of coal particles from gangue minerals by froth flotation occurs, for example, as dispersed air bubbles pass through a suspension of coal particles (-28 mesh). The bubble/particle aggregates of coal float to the surface and may be collected as a clean coal concentrate separated thereby from the wetted gangue particles.\nGenerally this process involves the use of suitable reagents (neutral molecular oils) to enhance the hydrophobic character of coal particles while the gangue mineral particles remain hydrophilic. These neutral molecular oils such as kerosene or fuel oil are called promoters and are used to enhance the attachment of air bubbles at the coal surface by forming a thin oil coating over the coal particles to be floated. Further, a frother is added to establish a stable froth phase to hold the bubble/particle aggregate. In coal flotation, frothers such as methyl isobutyl carbinol, terpinol, cresols, and polyglycols are frequently used. The choice of frother and oil depends on the desired level of selectivity with respect to ash and sulfur.\nAlthough the foregoing known processes are successful for conventional coal flotation, it is more difficult to produce \"super-clean\" coal by flotation of finely ground coal (-400 mesh). The promoter which is used to increase hydrophobicity inadvertently adsorbs on the gangue minerals, and these gangue minerals subsequently migrate to the concentrate, thus decreasing the quality of the product. Also, reagent consumption is high because fine coal, due to its high surface area, adsorbs significant amounts of promoter and frother. Finally, the rate of fine coal flotation by existing flotation techniques is very slow. As a result, the production of super-clean coal by conventional froth flotation methods has had limited success.\nA super-clean coal product is particularly desirable in the production of coal/water fuel. Coal/water fuel contains roughly 70% of the super-clean coal and is stabilized by the addition of various chemical additives so that it can be pumped, stored and used much like oil for which it is intended as a substitute.\nIn regard to the utilization of carbon dioxide in coal cleaning processes, some work has been done. For example, U.S. Pat. No. 4,522,628 to Savins discloses a method for removing ash from coal using liquid carbon dioxide under pressure in order to fracture and crush coal, not for flotation. In Savins, after comminution at high pressure and elevated temperature with liquid carbon dioxide, conventional flotation is used for coal recovery. Santhanam, in U.S. Pat. No. 4,206,610, also uses liquid carbon dioxide. The carbon dioxide is used in the Santhanam reference as a liquid to replace water as a medium to transport coal from mine to a remote processing plant. However, liquid carbon dioxide processes (transportation and cleaning) have inherent problems which relate to the chemical and physical properties of carbon dioxide (the need to keep it under positive pressure, etc.).\nOther references treating coal for various purposes include U.S. Pat. No. 3,998,604 to Hinkley which discloses the use of acids for grinding of coal. Hinkley uses acid treatment (not gas) during grinding. This is a type of leaching reaction. Carbon dioxide is mentioned in passing as a companion to carbonic acid for the sole purpose of grinding and making an acid slurry. Subsequently, the ground coal is floated and whether the coal \"floats\" or \"sinks\" is dependent on the flotation reagents, not the gas used.\nSteam and gaseous carbon dioxide have also been used in a high temperature coal process. Robinson et al, in U.S. Pat. No. 4,053,285, discloses a steam/carbon dioxide process for chemically treating (using hydrogen) previously prepared coal to lower the sulfur content. At high temperature and pressure, sulfur in the coal reacts with carbon dioxide and steam and is reduced to a sulfide. Flotation is not involved in this prior art reference.\nNo process is known which utilizes gaseous carbon dioxide to improve the hydrophobic character of the coal surface. That is, no process is known which uses gaseous carbon dioxide as a surface active reagent in the froth flotation cleaning of coal."}
{"text": "Computer software applications allow users to create a variety of documents to assist them in work, education, and leisure. For example, popular word processing applications allow users to create letters, articles, books, memoranda, and the like. Spreadsheet applications allow users to store, manipulate, print, and display a variety of alphanumeric data. Such applications have a number of well-known strengths including rich editing, formatting, printing, calculation, and on-line and off-line editing.\nUnfortunately for users, documents received by users may contain unnecessary or unwanted executable code embedded in the document. For example, a possible problem for computer and computer software users is receiving a document containing a “virus” in the form of an embedded executable code that executes when the user opens the document or performs some action in the document and which may cause harm to the user's document or to the user's computer software applications or computer hardware, or result in otherwise undesirable behavior. In word processor documents represented using an Extensible Markup Language (XML)-based file format, executable code can be located in various places, and finding such executable code in a fast and efficient manner becomes challenging. Due to its flexibility, XML has the ability to represent the same data in a virtually infinite number of ways. Accordingly, XML data representing executable code embedded in a document may be defined or represented in a number of different ways which makes locating the executable code difficult and time consuming. For example, XML supports the ability to encode text so that characters do not appear in their literal form, but which must be converted according to certain rules or according to an entity definition possibly existing elsewhere in the file. The definition for converting the coded text itself possibly refers to other components or entities with individual definitions existing in yet other places within the file or even existing in locations remote from the document or file.\nAdditionally, executable code may be placed almost anywhere in the XML file provided that it follows the rules associated with the XML structure of the file. This means that in order to find the executable code, a parsing or consuming application must first parse all the elements of the file prior to the embedded executable code. Ultimately, such a process may find the executable code within the XML formatted file, but the performance of the process is very slow, if not unacceptable, in environments where speed is critical, particularly in the case of virus checking prior to application or document startup, for example when the file is processed by an email server or an Internet gateway.\nIt is with respect to these and other considerations that the present invention has been made."}
{"text": "The present disclosure relates in general to electro-optical interfaces for transmitting and receiving optical signals through fiber optic cable systems. In particular, apparatuses and methods of assembling electro-optical interfaces are described that use a photonic integrated circuit for the active part of the printed circuit board to support optoelectronic components including vertical-cavity surface-emitting lasers (VCSELs) and photodiodes.\nAs technology progresses and the demand for high bandwidth transmission of optical signals over fiber optic networks in data centers increases, the tolerances of components in the optical interconnects become tighter, making optical interconnects more complex and more expensive."}
{"text": "Silicon integrated circuits continue to grow ever denser, requiring that the gate oxides in FETs (field effect transistors) be thinner and, in particular, that the gate electrodes in the FETs become narrower. Said gate electrodes are most usually made of polysilicon and, until relatively recently, were patterned by conventional use of photoresist masking.\nHowever, as gates have grown narrower, use of a photoresist mask has become increasingly problematical. There are a number of reasons for this--undercutting of the photoresist during etching, which can be tolerated in wider lines, introduces a degree of uncertainty into the process which can no longer be tolerated in very narrow lines; because of the high level of resolution required, the photoresist layer must be thinner than that used for etching wider lines; sensitivity of photoresist-based processes to process parameters such as exposure wavelength, brand of resist used, etc. becomes more of a problem; and line width control, and reduction of line edge roughness, through use of photoresist trimming, becomes more difficult if the photoresist is also used as the primary etch mask.\nAnother factor that needs to be considered when forming structures through etching is the need to use an ARC (anti-reflection coating). This is laid down just below the photoresist layer to prevent formation of standing wave patterns within the resist. ARCs may be either organic or inorganic, the former being relatively thick while the latter are relatively thin. For reasons that will become apparent, an inorganic ARC is part of the present invention.\nA routine search of the prior art was performed but no references that describe the exact process taught by the present invention were discovered. Several references of interest were, however, encountered. For example, Bell (U.S. Pat. No. 5,767,018) determined that a process being used by him to etch a silicon oxynitride ARC caused pitting of the polysilicon gate structure. Polymer formation during etching was considered to be a large contributor to this problem so an etchant that did not form a polymer was selected. This allowed a passivating coating to form on the sidewalls as the polysilicon was etched (using HBr/Cl.sub.2 /He--O.sub.2). The final step was an anisotropic etch that cleaned the top surface of the gate without removing the protective layer from the sidewalls.\nKumar et al. (U.S. Pat. No. 5,851,926) teach use of a tungsten silicide hard mask together with an NF.sub.3 and Cl.sub.2 etch to achieve vertical edges for a polysilicon gate. Keller (U.S. Pat. No. 5,346,586) etches a polysilicon layer that has been coated with a silicide layer by first forming an oxide hard mask and then successively etching the suicide and polysilicon layers.\nMuller (U.S. Pat. No. 5,674,409) shows a photoresist trim process for forming lines while Shin et al. (U.S. pat. No. 5,914,276) describe a process for etching a polycide/polysilicon gate."}
{"text": "The present invention concerns an electric linear actuator of the type comprising an electric motor, a tubular body comprising a piston and connecting means between the motor and the piston (commonly known as a jack).\nElectric linear actuators used until now have been driven by an electric motor disposed outside the body of the linear actuator rotating a screwthreaded rod through a gear, the rod driving the piston of the linear actuator.\nLinear actuators of the above kind are bulky, permeable because they are not sealed and fragile because the force is transmitted to the piston by a screwthreaded rod. There is a risk of the rod buckling when the linear actuator is operating in compression with a long travel.\nSealing against dust, inclement weather and moisture is difficult because the linear actuator is in two parts."}
{"text": "1. Field of the Invention\nThis invention relates to an improved method and apparatus for fluidized texturizing yarn. In a further aspect, the invention relates to a fluidized bounce crimp texturizing method, and apparatus therefor, wherein the thickness or a minute accumulation of the yarn issuing from the fluidized bounce crimp texturizer (bounce crimper) is sensed and the rate at which the texturized yarn collected regulated to maintain the thickness or limit the accumulation of the texturized yarn within certain predetermined limits.\n2. The Prior Art\nSynthetic thermoplastic yarn materials are produced as a number of continuous, straight, smooth filaments. Such yarns have little bulk, and their utility in textile applications is thus rather limited.\nIn order to enhance the bulk and texture of synthetic yarns, a variety of crimping processes have been used in the past. One common technique which has been used for thermoplastic yarns is to bend the yarn filaments and heat the yarn while the filaments are in bent or crimped configurations.\nOne especially good texturizing technique, in terms of yielding a high-bulked yarn, is known as \"rebound\" or \"bounce crimping\".\nBounce crimping entails hurling yarn, propelled by a heated fluid through a jet, in a continuous stream-like flow against a foraminous surface upon which the yarn impinges and from which the yarn instantaneously rebounds or bounces. The impact of the yarn upon the foraminous surface axially buckles and crimps individual filaments of the yarn while the heated fluid passes through the foraminous surface. The texturized yarn progresses without tension and substantially by rebound inertia away from the crimping zone and, in the prior art process, is guided to a collection station where the yarn is heated and then cooled to heat-set the crimp prior to winding upon a storage spool.\nThermoplastic yarn texturized by the foregoing bounce crimping process possesses, inter alia, exceptional covering capability and a high degree of resiliency; note U.S. Pat. No. 3,686,848.\nThe basic process and apparatus for practicing the bounce crimp texturizing process is described by Clarkson in U.S. Pat. No. 3,665,567. Brielfy, the Clarkson process entails feeding a yarn through an elongate slender tube by a jet of steam and hurling the yarn longitudinally against a foraminous screen causing it to buckle in a random manner and bounce away in a random array. The yarn is thereby crimped or texturized and freely rebounds laterally through a passage from which it drops down to a receiver for heat-setting. The steam primarily passes through the foraminous screen and is collected.\nIn order to properly conduct bounce crimping, it is important that the yarn is rebounded from the screen and discharged from the bounce crimper under essentially no tension. The yarn cannot be pulled from the bounce crimper, since tension on the yarn at this point could pull the crimp from the yarn and could also cause the foraminous screen to be bypassed.\nVarious improvements in bounce crimp texturizing processes and apparatus are described in U.S. Pat. Nos. 3,859,696, 3,859,697, 3,879,819 and 3,887,971. All of the apparatus and processes described in these patents and in U.S. Pat. No. 3,665,567 are characterized by the use of a J-tube type yarn accumulator wherein the yarn is accumulated (piled) and heat-treated (heat-set and cooled). In U.S. Pat. No. 3,879,819, the J-tube is provided with a photocell light sensing means for maintaining a certain height (pile) of yarn in the J-tube by regulating the yarn wind-up speed in response to the sensing means.\nThe J-tube accumulator was used by the prior art to heat-set the crimp on the yarn and to ensure that the yarn rebounded from the bounce crimp screen in a tensionless state by permitting the yarn to free-fall into the accumulator from the bounce crimper. Subsequently, the J-tube was primarily used only for the second purpose.\nIn accumulating or piling the yarn, tangles were found to occur, resulting in localized pulling on the yarn as it was wound up, thus causing the crimp to be pulled out of random segments of the yarn and/or the yarn to break. The frequency of these breaks necessitated an increase in the number of operators required to operate or monitor a given number of texturizing machines and rethread the yarn when breaks occur.\nAlso, where low-denier yarns (e.g., about 500 denier or less) are used, the problem is magnified such that accumulator systems cannot be efficiently used. This magnification is believed caused by the fact that the lower the denier, the more loops or coils that are in contact with each other in the accumulator. Hence, the more contact, the more chance there is for filaments of the various loops to tangle with each other. This increased contact, coupled with the lower weight of the loops or coils, substantially increases the likelihood of the loops being pulled out of the accumulator, resulting in increased piling and tangles, etc."}
{"text": "The present invention relates to analytical chemistry, and in particular relates to sample preparation for molecular analysis.\nIn order to carry out molecular analysis (the task of identifying one or more compounds in a sample) of any product, a sample of the product must be in such a form that it can be easily analyzed by chromatography, spectroscopy, mass spectroscopy and/or nuclear magnetic resonance instrumentation\nBecause these analytical instruments require substantially pure isolated analytes, some intermediate slops, generally referred to as “sample preparation”, must be carried out to isolate the compounds of interest from the sample matrix in which they might be found and prepare them for analysis by instrumentation.\nThe task of identifying one or more compounds in a sample—presents an enormously larger set of possibilities and challenges related to sample preparation. The number of “naturally occurring” compounds (those produced by plants or animals) is immeasurably large, and the capabilities of modern organic and inorganic synthesis have generated—figuratively or literally—a similar number of synthetic compounds.\nThere is tremendous interest in identification or quantitative measurement for compounds of interest as it relates to industrial processes, and for environmental testing for contaminants in waste water, soil, and air.\nEven a small group of recognizable representative samples would include pesticides in food, other synthetic chemicals in food (antibiotics, hormones, steroids), synthetic compositions (benzene, toluene, refined hydrocarbons) in soil, and undesired compositions in everyday items (e.g., Bisphenol-A (“BPA”) in polycarbonate bottles and other plastic food packaging.\nIn general extraction has been a main form of sample preparation; i.e., drawing one or more compounds of interest from a sample by mixing the sample with a solvent into which the desired compound(s) will be extracted from the sample so that it can be measured by an analytical technique.\nFor several generations (and continuing to date), sample preparation in the form of extraction has been carried out by the well-understood Soxhlet method which was invented in the 19th century. In the Soxhlet technique, a single portion of solvent circulates repeatedly through a sample matrix until extraction is complete. To the extent, the Soxhlet method has an advantage, it allows an extraction to continue on its own accord for as long as the boiling flask is heated and the condenser is cold.\nThis method of extraction can take hours to completely extract the compounds of interest. Other concerns of safely from flammable solvents, hazardous waste and breakable glassware are significant drawbacks to this method.\nAnother commonly known extraction method uses ultra-sonication; i.e., the irradiation of a liquid sample with ultrasonic (>20 kHz) waves resulting in agitation. Although ultra-sonication has an advantage of speeding up the extraction process the disadvantages are that it is a labor intensive, manual process and uses large amounts of solvents.\nIn more recent years analytical scale microwave-assisted extraction (MAE) has been utilized. MAE uses microwave energy to heat solvents in contact with a sample in order to partition analytes from the sample matrix into the solvent. The main advantage of MAE is the ability to rapidly heat the sample solvent mixture. When using closed pressurized-vessels the extraction can be performed at elevated temperatures that accelerate the extraction of the compounds of interest from the sample matrix. MAE accelerates the extraction process, but has its disadvantages as well. In the microwave healing process typically a polar solvent is needed to provide dipole rotation and ionic conduction through reversals of dipoles and displacement of charged ions present in the solute and the solvent, limiting non-polar solvent use. MAE uses expensive, high-pressure vessels that do not provide a means of filtering the extract, and they must be cooled before pressure can be released.\nIn the 1990's automated apparatuses for the extract ion of analytes were developed. These apparatuses incorporated solvent extraction in pressurized cells under elevated temperatures and pressures and are referred to as “Pressurized Fluid Extraction” (“PFE”) or “Accelerated Solvent Extraction” (“ASE”). PFE has shown to be similar to Soxhlet extraction, except that, the solvents are at elevated temperatures where they exhibit high extraction properties. This procedure was first developed by Dionex (Richter DE et al., Anal Chem 1996, 68, 1033). One such PFE automated extraction system (Dionex ASE) is commercially available.\nPFE was initially used for environmental contaminants (EPA Method 3545, herbicides, pesticides, hydrocarbons) in soil, sediments and animal tissues but has expanded to use in foods, pharmaceutical products and other biological samples.\nPFE provides an efficient extraction, but still has not overcome the major bottlenecks associated with the many steps necessary to prepare a sample for analysis. PFE utilizes multiple-component cells and many steps. The cells are tightly packed with the sample and other packing material to eliminate any void areas in the cell, enhance separation, and avoid channeling. Preparing a cell for analysis can typically take 15 minutes. The cells are pre-pressurized at pressures up to 1500 psi and heated up to 200° C. prior to adding the solvent. Extraction is based on chromatographic principles to force the hot solvent through the column. Cycle times can take up to 20 minutes and the requirements of high pressure lead to secondary disadvantages with respect to cost and maintenance.\nNewer PSE or ASE techniques attempt to address some of these difficulties, but still require that the cells be tightly packed, adding to the complexity and overall time required for each extraction.\nSample preparation, although having developed over the years, nevertheless remains the major bottleneck in molecular analysis. Accordingly, although the Soxhlet, Ultrasonication, MAE and PFE techniques have their advantages, each remains relatively time-consuming. As a result, when multiple samples are required or desired to provide necessary or desired information, the time required to carry out any given extraction-based molecular preparation step reduces the number of samples that can be prepared in any given amount of time, thus reducing the amount of information available in any given time interval. To the extent that, measurements are helpful or necessary in a continuous process, this represents a longer gap between samples or before an anomalous or troublesome result can be identified.\nIn recent decades, advances in liquid chromatography have led to analogous uses of packed columns in a technique referred to as solid phase extraction (“SPE”). Originally, chromatography was used to separate fractions in mixed samples for analytical purposes, and indeed it still serves this purpose very very well.\nIn SPE, the chromatography technique is modified to extract an anal vie from a matrix. Nevertheless, SPE fundamentally remains liquid chromatography technique in which molecules spread out (travel at different speeds) within a column based on their polarity, the particle size and polarity of the packed column (stationary phase), the polarity of the flowing liquid (mobile phase), the size (length and diameter) of the column and specific factors such as “hold-up volume,” “linear velocity,” and “flow rate.” See, e.g.. Arsenault, J. C. 2012. Beginner's Guide to SPE. Milford Mass.: Waters Corporation. (Arsenault 2012).\nAlthough SPE is useful, it has limiting characteristics, some of which include the following factors. First, a proper description of SPE is “liquid-solid phase extraction” because the sample matrix that holds the analyte is almost always a liquid.\nSecond, because SPE is essentially a liquid chromatography technique, it requires either column packing steps or a new column for each test, along with a potential pre-swelling step depending upon the material selected or required for the stationary phase. SPE typically requires different methods and manipulative steps for different analytes.\nThird a more deliberate (slower) flow through the packed column tends to produce belter separation among the fractions. Thus, in a very real sense slower SPE is better than faster SPE.\nFinally, if an additional driving force (i.e., in addition to simple gravity flow) is required to move solvent through the SPE column, an external liquid or gas pump, or a centrifuge, or a vacuum pull must be incorporated, which in turn increases, to some lesser or greater amount, the complexity of the technique and any supporting systems.\nMore recently, a dispersive solid phase extraction (“dSPE”) method referred to as “Quechers” or “QuEChERS” (“quick-easy-cheap-effective-rugged-safe”) has become a standard for extraction preparation of molecular samples. Dispersive SPE addresses some of the disadvantages of SPE, but still requires an extraction step, the adjustment of pH with an appropriate ionic salt, is labor-intensive (even if advantageous compared to other methods), and requires two separate centrifuge steps.\nQuechers is in many ways less complex than Soxhlet extraction, but still requires a multi-step process. In the literature, this is sometimes called a “three step process”(e.g., Paragraph 0153 of U.S. Patent Application Publication No. 20160370357), but in reality Quechers requires at least the following: homogenization of the matrix that contains the analyte of interest; adding extraction solvent; band agitation; buffering; a second agitation step; a centrifuge separation step; decanting; dispersive solid phase extraction (“dSPE”) clean up; a second centrifuge separation step; and decanting the supernatant liquid following the centrifuge step.\nIn addition to the multi-step handling and transfer of the solvent, the sample, and the various mixtures, each of the centrifuge steps takes a recommended five minutes; so that the full Quechers sample preparation takes at least about 15-20 minutes.\nAccordingly, although the Soxhlet, SPE, and Quechers (dSPE) methods have their advantages, each remains relatively time-consuming. As a result, when multiple samples are required or desired to provide necessary or desired information, the time required to carry out any given extraction-based molecular preparation step reduces the number of samples that can be prepared in any given amount of time, thus reducing the amount of information available in any given time interval. To the extent that measurements are helpful or necessary in a continuous process, this represents a longer gap between samples or before an anomalous or troublesome result can be identified.\nIn summary, among other disadvantages current sample preparation techniques are slow, require a large number of separate steps, use excess solvent, are difficult to automate, and operate under high liquid pressure.\nAccordingly, a need continues to exist for efficient rapid extraction-based molecular preparation techniques."}
{"text": "Inflammation may be an acute or chronic, localized or systemic immune and/or vascular response to trauma or infection by microbes, such as bacterial or viruses. Inflammatory reactions typically destroy, dilute, or confine the injurious agent and the injured tissue in the subject. Inflammation is characterized, particularly in the acute form, by the classic signs of pain, heat, redness, swelling, and possibly loss of function. At a histological level, inflammation involves a complex series of events, including dilation of arterioles, capillaries, and venules, and an increased permeability and blood flow, exudation of fluids, including plasma proteins, and leukocyte migration into the area of inflammation, particularly with a localized reaction.\nTherapeutic treatments for inflammation include a wide array of pharmaceutical drugs administered intravenously, subcutaneously, topically, or orally, depending on the particular inflammatory condition and outcome sought. However, most of the anti-inflammatory treatments available today have considerable drawbacks, including severe reactions at injection site, increased susceptibility to infection, rash, or other side effects. Thus, there is a need for better anti-inflammatory therapeutics and treatment methods.\nThymic Stromal Lymphopoietin (TSLP).\nThymic stromal lymphopoietin (TSLP) is an IL-7-like cytokine that triggers dendritic cell-mediated Th2-type inflammatory responses and is considered as a master switch for allergic inflammation. TSLP is an integral growth factor to both B and T cell development and maturation. Particularly, murine TSLP supports B lymphopoieses and is required for B cell proliferation. Murine TSLP plays a crucial role in controlling the rearrangement of the T cell receptor-gamma (TCR.gamma.) locus and has a substantial stimulatory effect on thymocytes and mature T cells. See, for example, Friend et al., Exp. Hematol., 22:321-328, 1994; Ray et al., Eur. J. Immunol., 26:10-16, 1996; Candeias et al., Immunology Letters, 57:9-14, 1997.\nTSLP possesses cytokine activity similar to IL-7. For instance, TSLP can replace IL-7 in stimulating B cell proliferation responses (Friend et al., supra). Although TSLP and IL-7 mediate similar effects on target cells, they appear to have distinct signaling pathways and likely vary in their biologic response. For Example, although TSLP modulates the activity of STAT5, it fails to activate any Janus family tyrosine kinase members (Levin et. al., J. Immunol., 162:677-683, 1999).\nTSLP Effects on Dendritic Cells and TNF Production.\nAfter human TSLP and the human TSLP receptor were cloned in 2001, it was discovered that human TSLP potently activated immature CD11c+ myeloid dendritic cells (mDCs) (see, e.g., Reche et al., J. Immunol., 167:336-343, 2001 and Soumelis et al., Nat. Immunol., 3:673-680, 2002). Th2 cells are generally defined in immunology textbooks and literature as CD4+ T cells that produce IL-4, IL-5, IL-13, and IL-10. And Th1 cells such as CD4+ T cells produce IFN-γ and sometimes TNF. When TSLP-DCs are used to stimulate naive allogeneic CD4+ T cells in vitro, a unique type of Th2 cell is induced which produces the classical Th2 cytokines IL-4, IL-5, and IL-13, and large amounts of TNF, but little or no IL-10 or interferon-γ (Reche et al., Supra) (see also, e.g., Soumelis et al., Nat. Immunol., 3:673-680, 2002). TNF is not typically considered a Th2 cytokine. However, TNF is prominent in asthmatic airways and genotypes that correlate with increased TNF secretion are associated with an increased asthma risk. See Shah et al., Clin. Exp. Allergy., 25:1038-1044, 1995 and Moffatt, M. F. and Cookson, W. O., Hum. Mol. Genet., 6:551-554, 1997.\nTSLP induces human mDCs to express the TNF superfamily protein OX40L at both the mRNA and protein level (Ito et al., J. Exp. Med., 202:1213-1223). The expression of OX40L by TSLP-DCs is important for the elaboration of inflammatory Th2 cells. Thus, TSLP-activated DCs create a Th2-permissive microenvironment by up-regulating OX40L without inducing the production of Th1-polarizing cytokines. Id.\nTSLP Expression, Allergen-Specific Responses and Asthma.\nin Early studies have shown that TSLP mRNA was highly expressed by human primary skin keratinocytes, bronchial epithelial cells, smooth muscle cells, and lung fibroblasts (Soumelis et al., Nat. Immunol., 3:673-680, 2002). Because TSLP is expressed mainly in keratinocytes of the apical layers of the epidermis, this suggests that TSLP production is a feature of fully differentiated keratinocytes. TSLP expression in patients with atopic dermatitis was associated with Langerhans cell migration and activation in situ which suggests that TSLP may contribute directly to the activation of these cells which could subsequently migrate into the draining lymph nodes and prime allergen-specific responses. Id. In a more recent study, it was shown by in situ hybridization that TSLP expression was increased in asthmatic airways and correlated with both the expression of Th2-attracting chomokines and with disease severity which provided a link between TSLP and asthma (Ying et al., J. Immunol., 174:8183-8190, 2005).\nTSLP Receptor (TSLPR) and Allergy, Asthma.\nThe TSLP receptor (TSLPR) is approximately 50 kDa protein and has significant similarity to the common γ-chain. TSLPR is a novel type 1 cytokine receptor, which, combined with IL-7Rα (CD127), constitutes a TSLP receptor complex as described, for example, in Pandey et al., Nat. Immunol., 1:59-64, 2000. TSLPR has a tyrosine residue near its carboxyl terminus, which can associate with phosphorylated STAT5 and mediate multiple biological functions when engaged with TSLP (Isaksen et al., J. Immunol., 168:3288-3294, 2002).\nHuman TSLPR is expressed by monocytes and CD11c+ dendritic cells, and TSLP binding induces the expression of the TH2 cell-attracting chemokines CCL17 and CCL22. Furthermore, as stated above, the TSLPR-induced activation of dendritic cells indirectly results in the increased secretion of TH2 cytokines IL-4, -5 and -13, which may be necessary for the regulation of CD4+ T cell homeostasis. In mice, deficiency of TSLPR has no effect on lymphocyte numbers. However, a deficiency of TSLPR and common γ-chain results in fewer lymphocytes as compared to mice deficient in the common γ-chain alone. See Reche et al., J. Immunol., 167:336-343, 2001 and Soumelis et al., Nat. Immunol., 3:673-680, 2002.\nStudies have found that TSLP and the TSLPR play a critical role in the initiation of allergic diseases in mice. In one study, it was demonstrated that mice engineered to overexpress TSLP in the skin developed atopic dermatitis which is characterized by eczematous skin lesions containing inflammatory infiltrates, a dramatic increase in circulating Th2 cells and elevated serum IgE (Yoo et al., J. Exp. Med., 202:541-549, 2005). The study suggested that TSLP may directly activate DCs in mice. In another study, conducted by Li et al., the group confirmed that transgenic mice overexpressing TSLP in the skin developed atopic dermatitis which solidifies the link between TSLP and the development of atopic dermatitis.\nAnother set of studies demonstrated that TSLP is required for the initiation of allergic airway inflammation in mice in vivo. In one study, Zhou et al. demonstrated that lunch specific expression of a TSLP transgene induced allergic airway inflammation (asthma) which is characterized by massive infiltration of leukocytes (including Th2 cells), goblet cell hyperplasia, and subepithelial fibrosis, and increased serum IgE levels (Zhou et al., Nat. Immunol., 6:1047-1053, 2005). However, in contrast, mice lacking the TSLPR failed to develop asthma in response to inhaled antigens (Zhou et al., supra and Al-Shami et al., J. Exp. Med., 202:829-839, 2005). Thus, these studies together demonstrate that TSLP is required for the initiation of allergic airway inflammation in mice.\nFurther, in a study conducted by Yong-Jun et al., it was demonstrated that epithelial cell-derived TSLP triggers DC-mediated inflammatory Th2 responses in humans which suggest that TSLP represents a master switch of allergic inflammation at the epithelial cell-DC interface (Yong-Jun et al., J. Exp. Med., 203:269-273, 2006).\nIn a recent study, it was shown that modulation of DCs function by inhibiting TSLPR lessened the severity in mice (Liyun Shi et al., Clin. Immunol., 129:202-210, 2008). In another set of studies, it was demonstrated that TSLPR was not only expressed in DCs, but also on macrophages, mast cells, and CD4+ T cells (Rochman et al., J. Immunol., 178:6720-6724, 2007 and Omori M. and Ziegler S., J. Immunol., 178:1396-1404, 2007). In order to rule out the direct effects of TSLPR neutralization on CD4+ T cells or other effector cells in allergic inflammation, Liyun Shi et al. performed experiments wherein OVA-loaded DCs were in vitro treated with anti-TSLPR before adoptive transfer to the airways of naive mice. It has previously been found that OVA-DCs triggered strong eosinophilic airway inflammation and accompanied with massive production of Th2 cytokines such as IL-4 and IL-5 (Sung et al., J. Immunol., 166:1261-1271 and Lambrecht et al., J. Clin. Invest., 106:551-559, 2000). However, pretreating OVA-DCs with anti-TSLPR resulted in a significant reduction of eosinophils and lymphocyte infiltration as well as IL-4 and IL-5 levels, further illuminating the role that TSLPR plays in DC-primed allergic disease. This result also supports that blocking of TSLPR on DCs will aid in controlling airway inflammation (Liyun Shi et al., supra).\nThere has been a growing body of experiments implicating the role of TSLP/TSLPR in various physiological and pathological processes. Physiological roles of TSLP include modulating the immune system, particularly in stimulating B and T cell proliferation, development, and maturation. TSLP plays a vital role in the pathobiology of allergic asthma and local antibody mediated blockade of TSLP receptor function to alleviate allergic diseases. Thus, interplay between TSLP and TSLP receptor is believed to be important in many physiological disease processes such as: allergic inflammation, skin lesions of patients with atopic dermatitis or atopic eczema, and asthma."}
{"text": "(1) Field of the Invention\nThe present invention relates to an electromagnetic induction-heating type fixing device and an image forming apparatus including such a fixing device, and in particular to a technology for improving a speed for fixing toner images on recording sheets of a small size (hereinafter, referred to as “small-sized recording sheets”) in a fixing device that fixes toner images on recording sheets of various sizes.\n(2) Related Art\nIn recent years, in the field of an image forming apparatus, saving of energy has been increasingly required as a part of global warming countermeasures, and accordingly, an induction-heating type fixing device that realizes high energy efficiency has been attracting attention.\nAn induction-heating type fixing device passes recording sheets through a fixing nip portion that is formed by a heating roller heated by induction heating and a pressurizing roller pressed against the heating roller so as to fuse and fix toner images thereon, for example.\nIn the fixing device that fixes toner images on recording sheets of various sizes, not only an excitation coil but also a demagnetization coil is used. The excitation coil heats the heating roller by induction heating over a width of the maximum sheet-passing region of the fixing device. The demagnetization coil prevents overheating of a part of a fixing roller in a region where, while small-sized recording sheets are passing through in the sheet-passing region, these sheets do not pass through (hereinafter, referred to as “non sheet-passing region”).\nFIG. 9 is an appearance perspective view showing an excitation coil and demagnetization coils pertaining to a conventional art. As FIG. 9 shows, demagnetization coils 901 are provided at positions corresponding to non sheet-passing regions that are aligned with both ends of fed small-sized recording sheets in a width direction thereof.\nWhen small-sized recording sheets are passed through, a circuit of the fixing device for applying current is closed and accordingly current inducted by a magnetic flux generated by an excitation coil 902 flows through the demagnetization coils 901. This causes the demagnetization coils 901 to generate a reversed polarity magnetic flux, which cancels the magnetic flux generated by the excitation coil 902. On the other hand, when recording sheets of a large size are passed through, the circuit is opened and accordingly demagnetization is stopped.\nAn image forming apparatus has been always required to improve a printing speed. It is thus necessary to increase a process speed (hereinafter, referred to as “fixing speed”) of a fixing device, that is, the number of recording sheets on which fixing is performed in units of time. In order to improve the fixing speed, more heat is naturally required, and accordingly output of an excitation coil is required to be increased.\nHowever, there is a problem that, if output of the excitation coil is increased, overheating of a heating roller in the non sheet-passing region becomes too extreme for practical use. FIG. 10 shows a temperature of the non sheet-passing region when recording sheets of an A6T size (105 [mm]×148.5 [mm]) are passed through in an image forming apparatus that can fix recording sheets of up to an A3 size. A horizontal axis of FIG. 10 represents a position (distance from the center of a sheet-passing region) in a direction perpendicular to a direction in which the recording sheets pass through, and a longitudinal axis of FIG. 10 represents a temperature of a fixing roller. Besides, a solid line 2101 and a dashed line 2102 indicate temperature distributions in the cases where speeds of passing through the recording sheets are 75 [ppm] and 65 [ppm], respectively. Note that ppm (papers per minute) represents the number of recording sheets that are passed for one minute.\nAs FIG. 10 shows, in the sheet-passing region, the fixing roller is deprived of heat by the recording sheets and heated by an excitation coil by electromagnetic induction at the same time, and accordingly a fixing temperature of the sheet-passing region remains at an appropriate temperature. On the other hand, since the recording sheets do not perform cooling of the fixing roller in the non sheet-passing region, a temperature of the fixing roller in the non sheet-passing region becomes higher than in the sheet-passing region.\nAlso, as the dashed line 2102 shows, when the speed of passing through the recording sheets is 65 [ppm], the highest temperature does not exceed 240° C., which is a general heat resistant temperature of silicone rubber. However, as the solid line graph 2101 shows, when the speed is increased to 75 [ppm], it turns out that the temperature exceeds 240° C. in some positions."}
{"text": "This invention relates to a liposome having active targeting property and ensured stability in blood, which can be used in the diagnosis and/or treatment of diseases, particularly renal diseases, that accompany production of various types of proteoglycan in injured portions of tissues and/or organs.\nIn recent years, there have been increased studies on a drug delivery system (DDS) based on a so-called targeting technique in which a drug is efficiently distributed in an organ of interest. Liposomes are one of the most extensively studied means as such a DDS because of their ability to include, or trap, drugs therein.\nHowever, there are various problems to be resolved in putting these liposomes into practical use, particularly, the difficulties in escaping from the mechanism of living body for recognizing foreign bodies and in controlling intracorporeal kinetics. That is, liposomes cause aggregation in blood by their mutual reaction with various blood plasma proteins, including opsonin, and are captured by a reticuloendothelial system (RES) such as in the liver or spleen, so that it has been difficult to deliver liposomes selectively to the target tissues or cells.\nIn recent years, it became possible to prevent aggregation of liposomes in blood and avoid being captured by RES, by coating the surface of liposome with a hydrophilic polymer such as polyethylene glycol (PEG) and thereby preventing adsorption of various blood plasma proteins, including opsonin, to the liposome surface (U.S. Pat. No. 5,013,556, U.S. Pat. No. 5,676,971).\nThere are a number of studies based on these techniques with regard to the targeting at tissues of interest, such as tumor tissues, regions of enhanced vascular permeability, inflamed tissues, the liver, the brain and lymph tissues, but all of these cases are based on the effect of so-called passive targeting property, which is induced as a result of the improvement of the stability in blood by avoiding the capture by RES (Advanced Drug Delivery Reviews, 24 (1997), 337-334), so that concern has been directed toward the development of a liposome having the ability to bind to the tissues of interest, in which the targeting function is reinforced by so-called active targeting techniques.\nRegarding the active targeting techniques, methods for the modification of the liposome surface with target factors such as antibodies, antibody fragments, amino acids, peptides or saccharides have been studied, as well as techniques for making the liposome surface into cationic nature. Particularly, the cation formation can be cited as one of the desirable modification techniques, because it has the ability to deliver genes and the like into cells, as a gene introducing technique, and it improves accumulation of liposome on, for example, an injured portion of vascular endothelium (JP-A-7-89874; the term xe2x80x9cJP-Axe2x80x9d as used herein means an xe2x80x9cunexamined published Japanese patent applicationxe2x80x9d). However, strong aggregation and the like phenomena mediated by the binding of blood protein to liposome were observed by the cationic liposome when compared with a neutral or anionic liposome (Biochimica et Biophysica Acta, 1280 (1996), 149-154), so that its sufficient active targeting property in the living body was not obtained as such.\nOn the other hand, proteoglycans which keep the cell surface anionic are known as the components which their interactions to cationic liposome are expected. It has been reported that overproduction of various types of proteoglycans occurs in tissues of fibrosis in a large number of organs (the liver, the lungs, the heart, the pancreas, the bone marrow and arteries) of tumors such as large bowel cancer, of cell proliferative nephritis and of other inflammatory diseases (Acta Pathol. Jpn., 36 (6): 827 (1986), FEBS Lett., 244: 315 (1988), J. Rheumatol., 18 (10): 1466, (1991), J. Dent. Res., 71 (9): 1587, (1992). This reaction in the living body is considered to be a result of excess repairing reaction which occurs during the step of wound healing or of excess cell proliferation in a tumor or the like tissue, so that it is used as a pathological marker of the foci of the aforementioned diseases.\nHowever, very little is known about reports which aim at achieving active targeting for tissues and/or organs in which such proteoglycan is overproduced. For example, JP-A-8-27030 describes that a drug carrier whose surface is modified with a basic compound which takes positive charge within a physiological pH range is accumulated on a cortical portion of the kidney, but there are no reports which disclose a targeting technique for the specific accumulation of a carrier on renal glomerulus, particularly renal glomerulus which caused inflammation by an injury, and proteoglycan producing tissues and/or organs typified thereby.\nGlomerulonephritis is one of the diseases so far reported in which production of various types of proteoglycan is accelerated in injured portions (Clin. Exp. Immunol., 108: 69, (1997), Kidney International, 49: supple. 54, s 55 (1996), J. Am. Soc. Nephrol., 2: s 88 (1992)). Importance of the targeting of a drug for a tissue and/or organ overproducing proteoglycan can easily be understood by merely citing this case of glomerulonephritis.\nAccording to the materials produced in 1996 by The Japanese Society of Dialysis Medicine, it is said that about 27,000 patients are subjected to dialysis due to renal insufficiency, and about 50% of the cases is caused by chronic nephropathy, and about 30% thereof by diabetic nephropathy. Since the primary disease in both of the chronic nephropathy and diabetic nephropathy is glomerulonephritis, these patients will be released from the struggling against dialysis when an effective method for the treatment of glomerulonephritis is established. Up to now, however, no effective glomerulonephritis treating method has been found.\nFor example, a glomerulonephritis, called IgA nephritis, is a glomerulonephritis most frequently found in the Japanese people, and it is said that its patients are estimated to be close to 300,000 in Japan alone and it occupies around 40% of the glomerulonephritis in adults and close to 30% of that in children. As a result of a follow up study carried out in Japan for recent 20 years, it was found that about 40% of IgA nephritis cases, estimated to be 5,000 to 6,000 cases per year (1995), reached terminal renal insufficiency and were subjected to dialysis, so that it is considered that the development of a therapeutic method specific for IgA nephritis is a world-wide pressing need. However, close to 5,000 patients have been subjected to dialysis every year because of the absence of sufficient means for its treatment.\nIt has been reported that, as a result of 10 years of clinical efficacy evaluation on IgA nephritis, the steroid therapy which is now actively used in Japan for nephrotic syndrome and rapidly progressive glomerulonephritis, and its efficacy is said to be established, also has a possibility of keeping renal functions for a prolonged period of time by its 2 to 3 years of continuous use at an early stage of progressive IgA nephritis.\nOn the other hand, however, side effects become a serious problem when it is continuously administered for a certain period for such a chronic disease, because problems such as arteriosclerosis, osteoporosis and immunity reduction in adults, slow development in children and the like cannot be avoided. As a result, it becomes a problem that its sufficient efficacy cannot be obtained, because not only QOL (quality of life) of the patient is reduced but also its administration cannot be continued due to side effects and, particularly, its administration for a prolonged period of time is impossible.\nIn order to resolve these problems, it is necessary to think out a means for expressing its drug effect alone, by a dose which does not express the aforementioned side effects, and it is desirable for this purpose to design a system for efficiently delivering a drug to the glomerulus which caused inflammation by an injury.\nThe same requirement is also applied to other diseases in which proteoglycan is excessively produced in injured portions of tissues and/or organs, so that it is expected to develop a method which can exert sufficient therapeutic effect without inducing side effects, by actively targeting the organ of interest with a drug which shows a strong drug effect but cannot be used due to its toxicity.\nIn consequence, this problem is resolved by including a drug into a liposome which is efficiently accumulated on tissues and/or organs that accompany production of proteoglycan, typified by renal glomerulus which caused inflammation by an injury, and by delivering the drug efficiently to such tissues and allowing it to perform continuous action. This purpose is achieved by realizing a liposome having active targeting property and ensured stability in blood, which also has the ability to be accumulated on injured portions of tissues and/or organs that accompany production of various types of proteoglycan.\nAccordingly, this invention contemplates providing a liposome having active targeting property and ensured stability in blood, which can be used in the diagnosis and/or treatment of diseases, particularly renal diseases, that accompany production of various types of proteoglycan in injured portions of tissues and/or organs.\nIn the process of carrying out studies on factors which control intracorporeal kinetics of liposome, the present inventors have found surprisingly that, when a liposome is prepared by incorporating a basic compound which takes positive charge within a physiological pH range, within a certain range of mixing ratio, and combining a hydrophilic polymer to its surface in an amount controlled within a certain range, it shows such properties that not only its stability in blood is ensured, its capture by the liver and the like can be avoided and its retention in blood is improved, but also its ability to recognize proteoglycan in injured portions is sharply improved.\nIt was found also that, when its major particle size range is from 90 to 200 nm, its accumulation on and selectivity and persistency of the effects are improved for tissues and/or organs in the living body, where proteoglycan is produced in excess amount, particularly for renal glomeruli. In this case, the term xe2x80x9cmajor particle size rangexe2x80x9d means a range of particle size in which 70% or more of particles are included when the particle size distribution of liposomes is defined by scattering intensity distribution in the particle size measurement based on the laser light scattering methods.\nThe present inventors have accomplished this invention by conducting intensive studies based on the aforementioned new knowledge.\nAccordingly, the present invention provides a liposome in which a drug is included, which comprises (1) a basic compound which takes positive charge within a physiological pH range, (2) a lipid derivative of a hydrophilic polymer and (3) a lipid which constitutes the liposome, as its membrane constituting components, wherein their constituting ratios are from 1 to 20 mol % of (1) based on (3) and from 0.2 to 5 mol % of (2) based on the total of (1) and (3).\nMore preferred is a liposome which comprises (1) a basic compound which takes positive charge within a physiological pH range, (2) a lipid derivative of a hydrophilic polymer and (3) a lipid which constitutes the liposome, as its membrane constituting components, wherein their constituting ratios are from 5 to 15 mol % of (1) based on (3) and from 0.2 to 5 mol % of (2) based on the total of (1) and (3).\nThe aforementioned basic compound which takes positive charge within a physiological pH range is preferably a basic compound which has an amidino group, a basic compound which has two or more amino groups, a basic compound which has a piperidine ring or a basic compound which has a quaternary amine, more preferably any one of those represented by the following formulae 1 to 4. \nIn this connection, the compound of formula 1 is described in WO 97/42166, the compounds of formulae 2 and 3 are described in JP-A-9-263579 and the compound of formula 4 is described in Chem. Papers, 39 (1), 125-134 (1985), but these documents describe nothing about a liposome having active targeting property and ensured safety in blood, which can be used in the diagnosis and/or treatment of diseases, particularly renal diseases, that accompany production of various types of proteoglycan in injured portions of tissues and/or organs. \n(In the formula 1, A represents an aromatic ring, each of R1 and R2 represents an alkyl or alkenyl group having from 10 to 25 carbon atoms, wherein R1 and R2 may be the same or different from each other, each of X1 and X2 represents xe2x80x94Oxe2x80x94, xe2x80x94Sxe2x80x94, xe2x80x94COOxe2x80x94, xe2x80x94OCOxe2x80x94, xe2x80x94CONHxe2x80x94 or xe2x80x94NHCOxe2x80x94, wherein X1 and x2 may be the same or different from each other, m is 0 or 1 and n is 0 or an integer of from 1 to 6.) \n(In the formula 2, R3 represents hydrogen or an alkyl or alkenyl group having from 1 to 8 carbon atoms, R4 represents hydrogen or an alkyl or alkenyl group having from 1 to 8 carbon atoms, each of R5 and R6 represents hydrogen or an alkyl or alkenyl group having from 1 to 25 carbon atoms (excluding a case in which R5 and R6 are both hydrogen), wherein R5 and R6 may be the same or different from each other, X3 represents xe2x80x94Oxe2x80x94 or xe2x80x94Sxe2x80x94, p is 0 or 1 and q is 0 or an integer of from 1 to 10.) \n(In the formula 3, each of R7 and R8 represents an alkyl or alkenyl group having from 1 to 8 carbon atoms, wherein R7 and R8 may be the same or different from each other, R9 represents hydrogen or an alkyl or alkenyl group having from 1 to 8 carbon atoms, each of R10 and R11 represents hydrogen or an alkyl or alkenyl group having from 1 to 25 carbon atoms (excluding a case in which R10 and R11 are both hydrogen), wherein R10 and R11 may be the same or different from each other, X4 represents xe2x80x94Oxe2x80x94 or xe2x80x94Sxe2x80x94, r is 0 or 1 and s is 0 or an integer of from 1 to 10.) \nThe aforementioned lipid derivative of a hydrophilic polymer is preferably a lipid derivative of polyethylene glycol, more preferably, said lipid derivative of polyethylene glycol is a compound which contains a polyethylene glycol chain and a diacyl glycerol in one molecule.\nIt is desirable that the aforementioned hydrophilic polymer has a molecular weight of from 1,000 to 7,000.\nRegarding the aforementioned liposome, it is desirable that its major particle size range is from 90 to 200 nm.\nIt is desirable that the aforementioned drug to be included in the liposome is a drug for use in the diagnosis and/or treatment of renal diseases or a drug for use in the diagnosis and/or treatment of diseases which accompany overproduction of proteoglycan in injured portions of tissues and/or organs.\nIt is desirable that the aforementioned drug for use in the diagnosis and/or treatment of renal diseases is an adrenocortical steroid and/or a derivative thereof.\nIt is desirable to use the aforementioned liposome as a drug for the treatment of inflamed kidney, and it is desirable to produce a medicament for the treatment of inflamed kidney making use of such liposomes.\nThe invention also provides a method for the treatment of inflamed kidney, which comprises administering these liposomes to mammals including human."}
{"text": "1. Field of the Invention\nThe invention relates to an air filter for internal combustion engines having a primary air region and a secondary air region disposed within a housing. The primary air comes from air located outside the internal combustion engine while the secondary air comes from crankcase emissions flowing from a crankcase in the internal combustion engine to the air filter.\n2. Description of the Prior Art\nVarious filters for reducing crankcase emissions are known an the art. For example, U.S. Pat. No. 5,125,940 to Stanhope, et al. discloses a venturi element 40 disposed outside of the downstream end of the housing 16 of the in-line air filter apparatus 10. Ambient air travels through tubular air filter 34, outside-in, and to the air intake hose 12 via the concave venturi element.\nU.S. Pat. No. 5,562,746 to Raether discloses an air filter assembly 10 having a plurality of filters 32 mounted to side walls 17 via tube sheet structure 28. A concave venturi 70 is fixedly coupled to each filter 32 by means of the tube sheet structure 28 and the venturi flanges 88.\nU.S. Pat. No. 4,272,263 to Hancock discloses a gas cleaning arrangement having a venturi sheet 4 disposed within a bag housing 1. Venturi sheet 4 supports a plurality of tubular filter arrangements 7 where each filter 7 includes a convex venturi 24, and accordingly the venturi is contained within housing 1.\nU.S. Pat. No. 5,383,356 to Zurek et al. discloses a concave venturi 76 that is separately connected to the outside of the downstream end of housing 86 or may be integrated with the housing as a single unitary element.\nU.S. Pat. Nos. 4,713,097 to Grawi et al 5,106,397 to Jaroszcyk, et al. and 5,120,334 to Cooper each disclose an air filter having a venturi, integrated with a housing, where the venturi is not adjacent to an air filter element.\nU.S. Pat. Nos. 4,272,262 to Britt et al. 4,834,784 to Biganset, and 5,196,653 to Kiss disclose air filters with venturi elements.\nThe present invention is an improvement over the prior art, because none of the prior art shows an air filter device having a venturi element shaped to be contained completely within the body of the housing and positioned such that crankcase emissions flow from a secondary air region past the primary air region and out of the housing. In addition, the prior art does not disclose the venturi element described above that is fixedly attached to the downstream end of an air filter element."}
{"text": "In the ear (ITE) hearing aids should be as small as possible and inserted as deeply as possible into the ear canal of the user. This is a general wish from the users of hearing aids. Today the limits of miniaturization are reached in that both the battery lid and the microphone must be placed in the faceplate of the shell. The faceplate is the part of the shell facing the surrounding environment, when the hearing aid is placed within the ear canal. In prior art hearing aids the battery lid and the microphone entrance are placed side by side, and therefore both take up space. The object of the invention is to provide a hearing aid, which is as small as possible, and which is insertable deeper into the ear canal than prior art hearing aids.\nIn behind the ear style hearing aids it is also a wish from the user that they be as small as possible. A further object of the invention is to provide a hearing aid wherein the microphone or microphones are accommodated in a space saving fashion."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit having an output circuit for outputting a balanced signal.\n2. Description of the Related Art\nAn output level of an output circuit for transmitting a signal from an integrated circuit to another integrated circuit is previously defined according to a standard. General examples include PECL (Pseudo-Emitter Coupled Logic), and LADS (Low Voltage Differential Signaling), and recently, PCI-express (Peripheral Component Interconnect-), XAUI (10 Gigabit Attachment Unit Interface), Infini Band, and Serial-ATA. FIG. 4 shows specification of typical interface levels. As apparent from FIG. 4, these interface levels are not compatible with each other. For example, comparing the specification of the PECL interface level with that of the LVDS interface level in FIG. 4, the output level (VOH, VOL) of PECL is a voltage lowered from a power supply voltage by a certain value, while the output level (VOH, VOL) of LVDS is a voltage which is independent of variation in the power supply voltage.\nAccordingly, the output circuit is generally configured according to a distinct circuit format suitable for the interface level of each standard. Each of these interfaces has a characteristic which cannot be achieved by other interfaces such as low power consumption. Thus, to use different interfaces for different purposes, there are many system devices having different interface levels with the similar function. As a result, transmission and reception between different interface levels is required. Such examples include an electrical input/output interface of an optical transmitter module.\nGenerally, the PECL or LVDS interface has become the mainstream of the electrical input/output interface of an optical transmitter module and is widely used in ASSP (Application Specific Standard Product) in many ways. In order to convert a signal into the interface level between PECL and LVDS as DC-coupled interfaces, the level using an external termination resistance is generally used. Hereinafter, an example is shown.\nFIG. 1 shows a circuit configuration of a typical example of a level converting method. In FIG. 1, an output of an output circuit 40 in the LVDS interface is converted into the PECL interface level and outputted to a receiver 50. The LVDS interface output circuit 40 has N channel transistors 41 and 42 as a differential pair, a current source 43, load resistances 46 and 47 having a resistance value RL and a level controller 48. A differential signal (INA, INB) is supplied to gates of the N channel transistors 41 and 42 of the differential pair and a signal of LVDS level as shown in FIG. 4 is output from the output terminals OUTA and OUTB.\nThe output of the output circuit 40 is converted into the PECL interface level by a level converting circuit having resistances 51 to 53, and 55 to 57, and supplied to the receiver 50. The resistance 51 having a resistance value R1, the resistance 52 having a resistance value R2 and the resistance 53 having a resistance value R3 are serially connected between a power supply voltage VDD2 and a ground GND. The output terminal OUTB is connected to a connection node of the resistance 52 and the resistance 53. A signal of the PECL level is outputted from a connection node ROUTB of the resistance 51 and the resistance 52. Symmetrically, the resistance 55 having the resistance value R1, the resistance 56 having the resistance value R2 and the resistance 57 having the resistance value R3 are serially connected between the power supply voltage VDD2 and the ground GND. The output terminal OUTA is connected to a connection node of the resistance 56 and the resistance 57. A signal of the PECL level is outputted from a connection node ROUTA of the resistance 55 and the resistance 56.\nGiven that the “H” level output voltage of the output nodes ROUTA and ROUTB is VOH, the “L” level output voltage of the output nodes ROUTA and ROUTB is VOL and amplitude, that is, differential output voltage of the output signal is VOD, each of the voltages can be obtained according to the following equations (1-1) to (1-3). With the power supply voltage VDD, VDD1=VDD2=VDD.VOH=VDD×(R2+R3)/(R1+R2+R3)+RL×I1×R1/{2×(R1+R2)}  (1-1)VOL=VDD×(R2+R3)/(R1+R2+R3)−RL×I1×R1/{2×(R1+R2)}  (1-2)VOD=RL×I1×R1/(R1+R2)  (1-3)\nBy properly selecting the resistances 51 to 53, 55 to 57, the level can be converted to correspond to the PECL interface to some extent. However, as understood from FIG. 4, the signal of the LVDS interface is a signal having a common voltage of 1.2 V and a fixed voltage independently from the power supply voltage. On the contrary, the signal of the PECL interface is a signal having a relative voltage which varies in level in connection with the power supply voltage. In the resistance-dividing level converting circuit in FIG. 1, as represented by Equations (1-1) and (1-2), the output voltages VOH and VOL vary according to a resistance division ratio of the power supply voltage VDD. Thus, the output voltages VOH and VOL satisfy the amplitude standard (VOD) of the PECL interface, but cannot satisfy the standard of the output level (VOH, VOL) unless the resistance values R1, R2 and R3 are changed depending on the power supply voltage VDD.\nFIG. 2 shows an example of a circuit for converting the PECL interface into the LVDS interface. An output circuit 60 of the PECL interface has transistors 61 and 62 as a differential pair, a current source 63, output transistors 65 and 66 and load resistances 67 and 68. A differential signal (INA, INB) is supplied to bases of transistors 61 and 62 and the signal of the PECL level in FIG. 4 is output from the output terminals OUTA and OUTB.\nThe output of the output circuit 60 is converted into the output of the LVDS level by a level converting circuit having resistances 71, 72, 74 and 75 and the converted output is supplied to a receiver 70. The resistance 71 having a resistance value R1 and the resistance 72 having a resistance value R2 are serially connected between the output terminal OUTA and the ground GND. The signal of the LVDS level is outputted from a connection node ROUTA of the resistance 71 and the resistance 72. Symmetrically, the resistance 74 having the resistance value R1 and the resistance 75 having the resistance value R2 are serially connected between the output terminal OUTB and the ground GND. The signal of the LVDS level is outputted from a connection node ROUTB of the resistance 74 and the resistance 75.\nGiven that a common voltage of balanced signals outputted from the output nodes ROUTA and ROUTB is VCM and an amplitude, that is, differential output voltage of the output signal is VOD, each of the voltages can be obtained according to the following equations (2-1) and (2-2):VCM=(VCC1−RL×I1/2−VF)×R2/(R1+R2)  (2-1)VOD=RL×I1×R2/(R1+R2)  (2-2)\nHere, VF is a base-emitter voltage of the transistors 65 and 66. In this example, a conversion reverse to the level conversion described referring to FIG. 1 is performed. By properly selecting the resistance values R1 and R2 according to equation (2-2), LVDS amplitude standard (VOD) can be satisfied. However, as represented by the equation (2-1), in accordance with change in the power supply voltage VCC1, the output common voltage VCM varies depending on the resistance division ratio. Thus, the level converting circuit cannot satisfy the standard of the output common voltage VCM unless the resistance values R1 and R2 changes in accordance with the change in the power supply voltage VCC1.\nSome interfaces do not use any external termination resistance. For example, as shown in FIG. 3, a typical example is a PCI-express interface. An output circuit 80 of the PCI-express interface has N channel transistors 81 and 82 as a differential pair, a current source 83 and load resistances 86, 87 and 88. A differential signal (INA, INB) is supplied to gates of the N channel transistors 81 and 82 and a signal of the PCI-express interface level is output from a connection node (OUTA, OUTB) of drains of the N channel transistors 81 and 82 and the load resistances 86 and 87, respectively. Outputs of the output circuit 80 are terminated by the termination resistance 91 having a resistance value RE and supplied to a receiver 90.\nThe PCI-express interface standard defines only output amplitude (VOD). When the output circuit 80 is used for the receiver 90 of the PECL interface, the resistance 88 having a resistance value RD may be adjusted to correspond to the “H” level output voltage VOH and the “L” level output voltage VOL of the PECL interface. The output levels VOH and VOL of the output terminals OUTA and OUTB and the amplitude VOD can be obtained according to the following equations (3-1) to (3-3):VOH=VDD1−{RL×RL/(2×RL+RE)+RD}×I1  (3-1)VOL=VDD1−{RL×(RL+RE)/(2×RL+RE)+RD}×I1  (3-2)VOD=I1×RL×RE/(2×RL+RE)  (3-3)\nThe equation (3-3) has a solution which satisfies the amplitude standard (VOD) of the PECL interface and the LVDS interface independently from the power supply voltage. As represented by the equation (3-1) and the equation (3-2), the output levels VOH and VOL of the output terminals OUTA and OUTB are determined based on the resistances 86 to 88 of the output circuit 80 and the terminal resistance 91 (resistance value RE) of the receiver 90. When the resistances 86 to 88 of the output circuit 80 are manufactured in the semiconductor integrated circuit together with the transistors and the like, the resistance values of the resistances have relatively a large manufacturing variation. Generally, it has been said that the resistance value of the resistance in the semiconductor integrated circuit has manufacturing variation of about −20% to +20%. Accordingly, when there is a mismatch between the above-mentioned resistance value and the resistance value of the termination resistance 91 on the receiving side, standards of the PECL interface, output levels VOH and VOL of the LVDS interface and the common voltage VCM cannot be satisfied. For example, to satisfy the standard of the output levels VOH and VOL of the PECL interface, the manufacturing variation of the resistance value needs to fall between −10% and +10%. Thus, it is difficult that the output circuit using the PCI-express interface satisfies the standard of the output levels VOH and VOL of the PECL or LVDS interface, or the standard of the common voltage VCM.\nAs described above, a load resistance (resistance value RL) in the output circuit and a current source (current value I1) have conflicting characteristics in variables (variation). For example, when the resistance value RL of the load resistance increases by 1.2 times due to the manufacturing variation, the current value I1 of the current source decreases by 1/1.2 times conversely. Accordingly, when the differential output terminals (OUTA and OUTB) are in the opened state, that is, nothing is connected to the output terminals, amplitude generated by the load resistance and the current source is kept constant in both the above-mentioned case.\nHowever, the output level standard of the PECL interface is linked to the power supply voltage and the output level standard of the LVDS interface is fixed with respect to the ground voltage, which have conflicting characteristics. When the level of the output of such an output circuit is converted by a level shift circuit having an external resistance inserted between the current source and the ground, the output level is determined with respect to the power supply voltage depending on the resistance division ratio. For this reason, the standard cannot be satisfied unless the resistance value is adjusted for each interface and each power supply voltage according to use environment.\nWhen a resistance built in the semiconductor integrated circuit is used as the load resistance of the output circuit, the resistance value greatly varies due to manufacturing variation. The output level of the output circuit is determined depending on a deviation ratio of the load resistance and a terminating resistance on the receiving side. For this reason, when the load resistance varies due to manufacturing variation and does not match the value of the terminating resistance on the receiving side, the standard of the interface cannot be satisfied. Especially in the PECL interface having a narrow allowable range of the output level, it is difficult to satisfy the standard.\nFIG. 5 summarizes compatibility of level conversion in the above-mentioned typical conventional output circuit. In FIG. 5, a circle represents availability including realization through level shift by the external resistance or the like or switching of current flow.\nJapanese Laid Open Patent Publication (JP-P2003-152522A) discloses a circuit for switching between PECL and LVDS, as a method which does not use the above-mentioned level conversion. An output circuit disclosed in Japanese Laid Open Patent Publication (JP-P2003-152522A) has a first output block including a first output port and a second output block including a second output port. The first and second output blocks are configured match a first transmission mode according to a first external control signal and bring about a first output characteristic in the first and second output ports. The first and second output blocks are configured match a second transmission mode according to a second external control signal and bring about a second output characteristic in the first and second output ports. The first transmission mode is a positive ECL (PECL) standard and the second transmission mode is a low-voltage differential signal transmission (LVDS) standard. Each of the first and second output blocks includes a switchable current source for feeding a current selected from a plurality of predetermined currents in the respective port according to the selected external control signal."}
{"text": "1. Technical Field\nThis application generally relates to generation of unique identifiers.\n2. Description of Related Art\nComputer systems may include different resources used by one or more host processors. Resources and host processors in a computer system may be interconnected by one or more communication connections. These resources may include, for example, data storage devices such as those included in the data storage systems manufactured by EMC Corporation. These data storage systems may be coupled to one or more host processors and provide storage services to each host processor. Multiple data storage systems from one or more different vendors may be connected and may provide common data storage for one or more host processors in a computer system.\nA host processor may perform a variety of data processing tasks and operations using the data storage system. For example, a host processor may perform basic system I/O operations in connection with data requests, such as data read and write operations.\nHost processor systems may store and retrieve data using a storage system containing a plurality of host interface units, disk drives, and disk interface units. Such storage systems are provided, for example, by EMC Corporation of Hopkinton, Mass. The host systems access the storage system through a plurality of channels provided therewith. Host systems provide data and access control information through the channels to the storage system and the storage system provides data to the host systems also through the channels. The host systems do not address the disk drives of the storage system directly, but rather, access what appears to the host systems as a plurality of logical disk units, logical devices, or logical volumes (LVs). The logical disk units may or may not correspond to the actual disk drives. Allowing multiple host systems to access the storage system allows the host systems to share data stored therein.\nIn connection with data storage system, and more generally any suitable system such as a computer system, it may be desirable to utilize techniques for generating unique identifiers for any one of a variety of different purposes."}
{"text": "1. Field of the Invention\nThe present invention relates to a thermal recording apparatus which thermally transfers inks of different colors on a print medium by selectively driving a plurality of heating-elements mounted on a thermal head to generate heat. The invention more particularly relates to a thermal recording apparatus in which odd and even heating-elements are energized at different time points to transfer an ink of one color every time each of the odd and even heating-elements are alternately energized, in order to prevent overlap of the first color ink and the second color ink, and then the third color ink is transferred on the position displaced from the positions where the first and second color inks are transferred, so that the color inks are less overlapped each other to improve color-reproducibility, and the size of apparatus and the manufacturing cost are reduced.\n2. Description of Related Art\nThere have been proposed various thermal recording apparatus to dissolve the problems such as the unevenness of density caused by a difference of transferring efficiency between inks and the mixing failure and fixing failure due to insufficient melted ink. Such the problems are caused when inks of different three colors, namely, cyan, magenta, and yellow are thermally transferred on a print medium by a thermal head so that heating-elements of the thermal head are energized by the same energy to thermally transfer the three color inks, overlapping one another.\nFor example, Japanese patent (JP) publication No. 3(1991)-54633 discloses a thermal transfer recording apparatus that is provided with a circuit for discriminating the recording history of a recording paper with a recorded picture element as a unit and a thermal head is controlled in accordance with the recording history to energize heating-elements to generate heat. The energy control is executed by the control of pulse width, duration of the power to be supplied to the heating-elements or voltage of the power, and the control of temperature of the thermal head. Accordingly, recording images transferred with uniform density can be obtained by the second and following thermal transfer recording, regardless of whether or not an ink is previously thermally transferred on the recording paper.\nJP patent publication, No. 3(1991)-73471 discloses a multicolor thermal recording apparatus provided with means for setting the energy to be supplied to a thermal head for the n-th overlap recording to be larger than that for the (nxe2x88x921)th recording. Concretely, the means sets the amount of energy to be supplied to the. thermal head for second overlap recording to be larger than that for the first recording, and the amount of energy for the third overlap recording to be larger than that for the second recording. Accordingly, the ink used for overlap recording can be melted in and mixed with the ink used before the overlap recording, improving mixability of the inks. Also, the ink used for the overlap recording can be saturated in the recording paper with the ink used before the overlap recording, so that the ink fixation can be enhanced.\nHowever, the above thermal recording apparatus disclosed in JP publications Nos. 3-54633 and 3-73471 require to control the amount of energy to be supplied to each of the heating-elements of the thermal head every time the thermal transfer is performed, and to change the amount of energy according to temperature variations in an operating environment. This results in a complicated control circuit, a large-sized apparatus and an increase in manufacturing cost.\nThe present invention has been made in view of the above-circumstances and has an object to overcome the above problems and to provide a thermal recording apparatus capable of reducing the amount of overlapping of plural inks having different colors by simple control, of improving color-reproducibility, and of reducing the apparatus size and the manufacturing cost thereof by energizing odd and even heating-elements at different time points to transfer an ink of one color every time each of the odd and even heating-elements are energized so that the first and second color inks are transferred without overlapping each other, and then the third color ink is transferred in a position displaced from the positions of the first and-second color ink.\nAdditional objects and advantages of the invention will be set forth in part in the description which follows and in part will be obvious from the description, or may be learned by practice of the invention. The objects and advantages of the invention may be realized and attained by means of the instrumentalities and combinations particularly pointed out in the appended claims.\nTo achieve the purpose of the invention, there is provided a thermal recording apparatus including a thermal head provided with a plurality of heating-elements, the thermal head being movable in a first direction, a drive device for selectively driving the heating-elements to generate heat, an ink ribbon applied thereon with at least a first color ink and a second color ink, a print medium on which images are to be printed through the ink ribbon by the selected heating-elements, the print medium being fed in a second direction intersecting the first direction, means for dividing the heating-elements into a fist group of odd heating-elements and a second group of even heating-elements, and alternately driving the odd heating-elements and the even heating-elements at different time points, a memory for storing image data corresponding to the first color and image data corresponding to the second color, and control means for operating the drive device to drive one group of the first heating-element group and the second heating-element group to generate heat and, after moving the thermal head by one pitch defined by a print dot in the first direction, operating the drive device again to drive the other group of the first heating-element group and the second heating-element group to generate heat, thereby to print the first color ink on the print medium based on the image data, and for operating the drive device to drive the other group that is different from the one group driven first to print the first color ink thereby to print the second color ink on positions where the first color ink has been printed on the print medium.\nIn the above thermal recording apparatus according to the present invention, the heating-elements associated with odd dots to be printed and those associated with even dots are driven to generate heat at different time points. To thermally transfer an ink of the first color on a print medium based on the image data, a plurality of heating-elements out of one of the odd and even heating-elements are energized and, after the thermal head is moved by a printing pitch in the predetermined sub-scanning direction, a plurality of heating-elements out of the other of the odd and even heating-elements are energized. In this way, the odd and even heating-elements are alternately energized every time after the thermal head is moved by a printing pitch each.\nSequentially, to thermally transfer an ink of the second color, the odd or even heating-elements that differ from ones used for the thermal transference of the first color ink in a transference start position thereof are energized to generate heat, transferring the ink of the second color.\nAccordingly, the first color ink and the second color ink are thermally transferred without overlapping each other, so that images with uniform density can be formed on the print medium. Since the energy to be supplied to a heating-element needs no controlling for respective heating-elements, each of the first and second color inks can be thermally transferred on the print medium by a simple control, enabling the reduction of the size and manufacturing cost of the apparatus. The first color ink and second color ink are thermally transferred in a staggered and latticed arrangement, thereby improving color-reproducibility due to an additive process. Furthermore, the alternate driving of the odd and even heating-elements can substantially prevent the influence of accumulation of heat on each of the heating-elements, which enables printing of color dots with a fine diameter."}
{"text": "One of the known problems associated with many industrial factories operating under elevated temperatures is the generation of excess heat emitted from the facility, which is not effectively utilized.\nAlthough there could be a substantial potential of energy content in the outgoing gases, still, the typical problem associated with utilizing this potential is the fact that the hot gaseous feed stream is a by-product generated on a transient basis, e.g., an outcome of batch processes, non-constant generation of the by-product hot gases, and the like. Thus, it would be virtually impossible to rely on an essentially non-constant hot gases feed or even intermittent feed of hot gases, for running a continuous process using this heat, which is desired from economical point of view.\nOne way of solving this problem, is by including a storage in the system for storing the hot gases leaving the process and drawing a constant hot gaseous feed from that storage. However, this solution is impractical due to the costs associated with the need to store and compress the gases at elevated temperatures.\nAnother way to obtain an intermediate storage is to use a phase changing material in which the heat can be stored. However, once again the costs involved with such a process and the loss of process efficiency (due to the reduction of the temperature in the storage and retrieval processes), render this solution impractical.\nYet another way to obtain such a storage is by using materials which have high heat capacity (e.g. refractory bricks through which hot gas carrying the excess heat flows) and by heating this media, the excess heat is stored as sensible heat. Again, this solution suffers from similar disadvantages as the solutions referred to above.\nIn view of that, there is a need to obtain a solution that would enable utilizing the excess heat and consequently to increase the overall efficiency of the process in which the excess heat is generated."}
{"text": "Free radical olefin polymerization is well-known and commercially important. However, copolymerizations of ethylene with polar industrial monomers, such as vinyl acetate and acrylic acid, are increasingly complex, requiring high temperature and exceedingly high pressures. Further, such processes intrinsically lack precise control over the resulting polymer architecture, polymer molecular weight and incorporation of polar monomer. In view of these deficiencies, coordination polymerization has been explored for potential controllable strategies for the synthesis of polyolefins having various functionalities derived from the incorporation of a polar monomer. Two dominant classes of transition metal catalysts have been developed to date for copolymerization of ethylene and industrial polar monomers. The first class encompasses group 10 complexes coordinated by an α-diimine ligand, commonly referred to Brookhart-type catalysts. The remaining class employs neutral palladium complexes coordinated by a phosphine sulfonate (Drent-type). These two classes have persistent limitations. For Brookhart catalysts, complex stability is limited for polymerizations conducted above room temperature. Even state-of-the-art Brookhart catalysts are persistent for only about 15 minutes at or above 90° C. While Drent catalysts generally exhibit greater thermal stability, they typically produce low-molecular weight copolymers of ethylene and polar industrial monomers and/or turnover with poor rates. Therefore, new transition metal catalysts are required for the production of polar functionalized polyolefins."}
{"text": "A variety of approaches to developing driverless vehicles focus predominately on automating conventional vehicles (e.g., manually-driven automotive vehicles) with an aim toward producing driverless vehicles for consumer purchase. For example, a number of automotive companies and affiliates are modifying conventional automobiles and control mechanisms, such as steering, to provide consumers with an ability to own a vehicle that may operate without a driver. In some approaches, a conventional driverless vehicle performs safety-critical driving functions in some conditions, but requires a driver to assume control (e.g., steering, etc.) should the vehicle controller fail to resolve certain issues that might jeopardize the safety of the occupants.\nAlthough functional, conventional driverless vehicles typically have a number of drawbacks. For example, a large number of driverless cars under development have evolved from vehicles requiring manual (i.e., human-controlled) steering and other like automotive functions. Therefore, a majority of driverless cars are based on a paradigm that a vehicle is to be designed to accommodate a licensed driver, for which a specific seat or location is reserved within the vehicle. As such, driverless vehicles are designed sub-optimally and generally forego opportunities to simplify vehicle design and conserve resources (e.g., reducing costs of producing a driverless vehicle). Other drawbacks are also present in conventional driverless vehicles.\nOther drawbacks are also present in conventional transportation services, which are not well-suited for managing, for example, inventory of vehicles effectively due to the common approaches of providing conventional transportation and ride-sharing services. In one conventional approach, passengers are required to access a mobile application to request transportation services via a centralized service that assigns a human driver and vehicle (e.g., under private ownership) to a passenger. With the use of differently-owned vehicles, maintenance of private vehicles and safety systems generally go unchecked. In another conventional approach, some entities enable ride-sharing for a group of vehicles by allowing drivers, who enroll as members, access to vehicles that are shared among the members. This approach is not well-suited to provide for convenient transportation services as drivers need to pick up and drop off shared vehicles at specific locations, which typically are rare and sparse in city environments, and require access to relatively expensive real estate (i.e., parking lots) at which to park ride-shared vehicles. In the above-described conventional approaches, the traditional vehicles used to provide transportation services are generally under-utilized, from an inventory perspective, as the vehicles are rendered immobile once a driver departs. Further, ride-sharing approaches (as well as individually-owned vehicle transportation services) generally are not well-suited to rebalance inventory to match demand of transportation services to accommodate usage and typical travel patterns. Note, too, that some conventionally-described vehicles having limited self-driving automation capabilities also are not well-suited to rebalance inventories as a human driver generally may be required. Examples of vehicles having limited self-driving automation capabilities are vehicles designated as Level 3 (“L3”) vehicles, according to the U.S. Department of Transportation's National Highway Traffic Safety Administration (“NHTSA”).\nAs another drawback, typical approaches to driverless vehicles are generally not well-suited to detect and navigate vehicles relative to interactions (e.g., social interactions) between a vehicle-in-travel and other drivers of vehicles or individuals. For example, some conventional approaches are not sufficiently able to identify pedestrians, cyclists, etc., and associated interactions, such as eye contact, gesturing, and the like, for purposes of addressing safety risks to occupants of a driverless vehicles, as well as drivers of other vehicles, pedestrians, etc.\nThus, what is needed is a solution for implementing autonomous vehicles, without the limitations of conventional techniques."}
{"text": "1. Field of the Invention\nThe present invention relates to a scintillator, a material that emits light when excited by radiation. More particularly, the present invention relates to a scintillator crystal functionally capable of guiding light emitted from the scintillator to photodetectors. The present invention also relates to a radiation detector using a scintillator crystal.\n2. Description of the Related Art\nIn a X-ray CT (computed tomography) apparatus to be employed in a clinical site, X-rays transmitted through a subject are received by scintillators and lights emitted from the scintillators are detected by photodetectors. The photodetectors are arranged in two-dimensional array and the scintillators are separated by partitions so as not to give rise to any crosstalk of light. The partitions are desirably formed as thin as possible from the viewpoint of not contributing the partitions to a detection of X-rays and preventing the spatial resolution of the X-ray CT apparatus from degrading.\nFor example, Japanese Patent Application Laid-Open No. 2008-145335 discloses a technique of forming a scintillator array by bonding together a large number of scintillator crystals by means of adhesive, subsequently etching the adhesive and filling the voids produced as a result of the etching with titanium oxide powder as partition material. According to the cited patent literature, the partitions as thin as about 1 μm can be formed.\nHowever, while the technique disclosed in Japanese Patent Application Laid-Open No. 2008-145335 can produce thin partitions, it cannot eliminate the presence of partitions. Also, from the viewpoint of manufacturing steps, the disclosed technique requires a large number of steps including a step of cutting a scintillator and a step of bonding scintillators to form the partitions.\nFurther features of the present invention will become apparent from the following description of exemplary embodiments with reference to the attached drawings."}
{"text": "In some instances, it is desirable to configure a system to receive all of the electromagnetic signals within the receiver's capability dictated by its sensitivity and bandwidth limitations. Usually the signals of Interest are incident from widely diverse directions. Thus, in one principal method, an antenna having a wide azimuth beamwidth, such as an omnidirectional antenna, is chosen as the system's receptor element.\nThe low-directive gain of such antennas is a severe limitation on the system sensitivity. In addition, the wide beamwidth does not permit directional resolution of multiple signals. Such resolution is usually desired to prevent garbling of signals that cannot otherwise be resolved in frequency or time-of-occurrence. Directional resolution is also useful in cases where the direction of incidence of the signals is to be estimated.\nTo overcome these disadvantages, alternative systems have been configured using high-gain narrow-beam antennas. In one such system, a narrow-beam is scanned over the azimuth coverage sector (either by mechanical motion of the antenna or by use of electronically controllable phase shifters). The disadvantage of this system is that the beam cannot look everywhere at once. This is especially a problem for multiple signals from diverse directions if they have rapidly changing waveforms (high information rate or short-pulse signals). These signals may hoe be sampled at a sufficiently high rate by the scanning beam to prevent information loss.\nIn another system, the antenna is configured to form multiple, narrow beams, that are contiguous and fixed in direction. Each beam port of the antenna is connected to a separate receiver. If the number of beams is sufficient to cover the entire azimuth sector of interest, then the system can exhibit the advantages of high directive gain (high sensitivity), good directional resolution (direction estimate and suppression of garble), and complete, simultaneous directional coverage (no information loss). However, the disadvantage of this system is the high cost of the multiple receivers.\nStill another group of alternative systems have been configured to overcome these disadvantages by scanning a narrow beam at scanning rates so high that the beam will intercept each signal at least twice as often as its information rate. Thus, in the case of pulsed signals, the beam is caused to scan through its complete coverage sector within the time period of the shortest pulse expected. Such rapid scan is obtained by heterodyne techniques. This results in bandwidth spreading of the received signals, which, in the case of a pulsed signal, yields a predictably compressed pulse whose time of occurrence is directly related to the emitter azimuth location.\nEarly embodiments of array antennas using this rapid-scan heterodyne technique suffered a sampling loss which degraded sensitivity. This loss of signal (relative to the maximum signal energy available) is caused by the fact that the signal received by each radiator is summed coherently for only the portion of the scanning period during which the beam is pointing nearly in the direction of signal incidence. During the rest of the scan period, the signal received by each radiator is coherently summed at the resistor-terminated ports of the combiner (if it is a resistively isolated type of combiner) or is reflected back from the combiner to the radiators (in the case of a non-isolated type of combiner). The sampling loss in dB for an N element array is given by 10 log N; this is approximately the same numerical value as the gain of the array (relative to the gain of a single element). Although this antenna system can resolve emitters in azimuth as well as a conventional N element array, it provides only the sensitivity that can be attained with only a single element of that array.\nA recent embodiment of an array antenna using this rapid-scan heterodyne technique (patent application Ser. No. 519,161, filed August 1983 by R. M. Rudish) overcame the sampling loss deficiency of the earlier embodiments by providing a multibeam capability. The earlier embodiments scan a single beam, whereas the recent embodiment uses a multiple beam-forming device such as a Butler matrix to produce and scan a comb of multiple contiguous beams in unison. At any instant, the beams are spread over the entire coverage sector of the antenna. At all times, at least one of the beams will be directed at any emitter within the antenna's coverage. The beams are displaced from each other in direction by approximately their beamwidths. They are also displaced from each other in time; that is, the time at which a particular beam coincides with a given direction differs from that at which each of the other beams coincide with that same direction. The time differences are known, fixed amounts related to the scan rate. Thus, by using a tapped delay line summing network, the recent embodiment can impart just the right differential delays to the individual beams to make their outputs from any one emitter occur in unison so that they can be summed coherently. In this manner, the recent embodiment recovers the signal loss which occurred in the arrays of earlier systems due to sampling the signal during only a portion of its time of presence. Thus, full array gain can be realized for an increase in system sensitivity. Yet, because there is only a single, summed output, far fewer components are needed to complete the processing of that output than is the case for other types of multiple-beam-antenna/multiple-receiver systems.\nThe recent embodiment of this rapid-scan heterodyne technique will be described in greater detail for the case of an N-element array since it forms the foundation for the improvements of the present invention."}
{"text": "Restaurants and institutions commonly hold and serve hot foods from electrically operated countertop food warmers. Water is usually, but not necessarily, used in the wells of the warmer to act as a heat transfer medium and to improve thermal efficiency. These warmers are designed to hold 12\".times.20\" steamtable pans or combinations of fractionally sized pans in various depths. With the use of an adaptor plate formed from a sheet of metal with one or more openings. These warmers may also be used to hold a variety of round-shouldered vegetable pans or inserts. The rectangular vegetable pans or inserts are supported by a lip which extends outwardly at the top of the pan and rests on the top edge of the warmer well. In addition to support, the lip acts as a loose seal to prevent the escape of large amounts of steam from the area between the bottom of the well and the bottom of the pan. Like the rectangular pans, the adaptor plates also have lips and flanges to provide support and act as a steam seal.\nIn existing food warmers, the arrangement is as follows. A food warming insert, pan or vessel occupies the innermost position for the purpose of holding food. Just outside of the insert or pan is a metal well in which the pan is supported, and it is outside this well that an electric heating element is located. An electric band heater is often secured against the outside surface of the well to keep the food contained in the vessel at the desired temperature. In other designs, an electric heating element is spaced a short distance below or to the side of the vessel. A housing containing insulation usually surrounds the heating element. A major shortcoming of this arrangement is the heat loss resulting from having to force the heat through the well, then through the pan into the food. A general objective of the present invention is to find a way of placing the electric heating element inside the well and for efficiently transferring heat to the food containment vessel or pan in three ways; namely, by conduction, convection and radiation to thereby improve thermal efficiency of the food warmer and save energy."}
{"text": "Since continuously porous foams made of polyurethane are light and have functions of absorbing liquids, storing liquids, permeating liquids and gases, absorbing sounds, selectively removing solids, absorbing shocks and the like, they are applied to ink rollers, self-inking rubber stamps, writing-brush pens, filters, sponges for cosmetics and the like.\nFoams made of polyurethane have been mainly produced by chemical foaming methods and the like. However, since, in the chemical foaming methods, pores are formed with a gas generated in chemical reaction, it has been impossible to make the so-called cell-diameters 200 μm or less on average and the resultant foams have been unable to effectively absorb fine water droplets. Therefore, continuously porous foams produced by extraction methods using a pore generation agent are proposed.\nFor example, JP52-32971A (Patent Literature 1) discloses a porous molding produced by kneading a composition consisting of a thermoplastic polyurethane, calcium carbonate in a form of powder soluble to an acid, a surfactant or polyhydric alcohol, dimethylformamide and acetone, thermally molding the kneaded composition at 130° C., and immersing the resultant molding in an aqueous hydrochloric acid solution and then washing and drying.\nIn this process, the calcium carbonate is added as a pore generation agent, and pores are formed by extracting the calcium carbonate dispersed in the molding with acid and water followed by washing. In addition, the surfactant is added to facilitate extraction of the calcium carbonate with water.\nAs an alternative, JP58-189242A (Patent Literature 2) discloses a porous polymer foam produced by a process of dissolving a polyurethane in a solvent such as dimethylformamide, mixing the resultant solution with a composition compounded with a pore generation agent such as polyvinyl alcohol, filling in a predetermined mold to coagulate the polyurethane in an non-solvent for the polyurethane, followed by extracting out the pore generation agent with a large quantity of water.\nCurrently, in manufacturing precise products such as printed circuit boards and lead frames, continuous pore elastomers made of polyurethane are increasingly demanded in applications to water-absorbing rollers for cleanly and uniformly draining off water droplets adhered on the surfaces of water-washed product and to swabs for absorbing water adhered on the precise products. The continuous pore elastomers used for such applications are required to have a property of instantly absorbing water. Furthermore, for the application to water-absorbing rollers, also required is a property effectively spitting out water in compressing the rollers with an external force. However, it has been difficult for the above-mentioned continuous pore elastomers to provide such properties and they have not suited to the applications such as water-absorbing rollers and swabs. Accordingly, the development of continuous pore elastomers having excellent properties such as instantly absorbing water has been desired.\nIn addition, since such continuous pore elastomers have become widely used for applications of daily life and industries, the development of methods for producing the continuous pore elastomers massively, inexpensively and stably is desired. The above-mentioned production process described in Patent Literature 1, however, requires a heating at a high temperature in a step of producing kneaded composition, resulting in production cost increase and causing ingredient degradations. Furthermore, use of an organic solvent with low-boiling point causes a possibility of taking fire at a high temperature. Moreover, use of the aqueous hydrochloric acid solution requires neutralizing a liquid wasted after the use, also causing a cost increase.\nOn the other hand, the process of Patent Literature 2 requires about 1 week until the coagulation completes, for example, when the molding has a thickness of 20 mm, and also needs a mold for molding fabricated with an expensive porous material in its production. Particularly, for mass production, a great number of such molds are required, resulting in increase of the production cost. Therefore, the development of a process for producing a continuous pore elastomer, which is free from the problems of the conventional techniques, is desired.\nPatent Literature 1: JP52-32971A\nPatent Literature 2: JP58-189242A"}
{"text": "Before discussing the invention in more detail, it is useful to consider some background. PNNI is a hierarchical, dynamic link-state routing protocol defined by the ATM Forum for use in ATM networks. The PNNI protocol provides, inter alia, a system for creation and distribution of topology information which determines how individual network nodes “see” the network and thus how nodes communicate. A key feature of the protocol is the ability to cluster groups of switches into “peer groups”. The details of each peer group are abstracted into a single logical node (a “logical group node” or LGN) which is all that can be seen outside of that peer group. One node in each peer group serves as the “peer group leader” and represents that peer group as the logical group node in the next level up of the hierarchy. This system is applied recursively so that PNNI can hierarchically aggregate network topology information. The PNNI topology information available to switches is such that each switch sees the details of its own peer group plus the details of any peer group that represents it at a higher level of the PNNI hierarchy. It is this hierarchical topology abstraction that reduces the resources required to support large-scale ATM networks.\nThe topology data communicated over PNNI networks is defined by PNNI Topology State Elements (PTSEs). PTSEs include data relating to nodes, links and addresses which can be accessed by network devices, and are created and distributed by nodes so that each node can maintain a topology database which defines its view of the network. PTSEs are flooded among nodes in a peer group so that each peer group node has the same topology database and thus the same view of the network. In the next level up of the hierarchy, however, the peer group topology is abstracted into a single logical node as described above. The logical group node generates PTSEs advertising addresses accessible within its child peer group and distributes these to its neighbors in the next level of the hierarchy, but the details of nodes and links within the peer group are lost. PTSEs generated by a logical group node in this way are also flooded back down the hierarchy, together with PTSEs received by the LGN from its neighbors, to enable the lower-level nodes to identify their “ancestors” (i.e. their representative nodes at higher levels of the PNNI hierarchy) and maintain their views of the PNNI topology.\nIn general, a PTSE which has been flooded through the network can only be modified by its originating node (i.e. the node which originated, or generated, that PTSE), though PNNI does define a “time-out” system for PTSEs whereby each PTSE is given a lifetime for which it is valid, normally one hour. A PTSE's originating node should “refresh” the PTSE periodically by redistributing the PTSE to its neighbors, so that the PTSE is again flooded through the network. However, if a PTSE's lifetime expires without the PTSE being refreshed, the PTSE is no longer considered valid topology information and is removed, or “flushed” from the network. So if a node becomes inaccessible due to a link failure for example, then PTSEs relating to that node will eventually be flushed from the network. In addition, our copending European Patent Application No. 99115580.5, filed 6 Aug. 1999, discloses mechanisms whereby a peer group leader can check the accessibility of addresses in its own peer group and notify neighboring nodes of changes in address accessibility within that peer group.\nPNNI provides full support for mobility at the ATM layer (“PNNI Addendum for Mobility Extensions v1.0”, ATM Forum af-ra-0123.000, April 1999). For example, the PNNI mobility extensions allow a logical group node abstracting a mobile ATM network to roam in the PNNI hierarchy of a terrestrial backbone network Routing information detailing the current location of the mobile network is advertised through regular PNNI, thus enabling the establishment of calls from a terrestrial end-system to an end-system of the mobile network, and vice versa. In addition, ATM networks can be used to carry higher layer protocol information such as IP (Internet Protocol) information. This can conveniently be done by employing an extension to the PNNI protocol known as PAR (PNNI Augmented Routing). PAR is described, for example in “PNNI Augmented Routing (PAR)”, af-ra-0104.000, ATM Forum, January 1999. Briefly, PAR allows IP information, which is not related to operation of the ATM network in itself, to be distributed over the network. PAR makes use of the PTSEs discussed above to distribute IP-related information in addition to the ATM topology information. PAR-enabled devices in the network encapsulate UP information in PTSEs which are then distributed in the usual PNNI way. The IP information in these so-called “PAR PTSEs” is opaque to PNNI nodes that are not PAR-enabled, but other PAR-enabled nodes are aware of the format of the IP information in PAR PTSEs. Thus, a PAR-enabled device in the network can communicate IP information over the network by means of PAR PTSEs, and another PAR-enabled device can extract the IP information.\nA further extension of the PNNI protocol known as “Proxy-PAR” allows higher layer protocol devices, in particular IP devices such as routers, to communicate IP information over the network without themselves participating in PNNI. Proxy-PAR is also described in “PNNI Augmented Routing (PAR)”, af-ra-0104.000, ATM Forum, January 1999. Briefly, Proxy-PAR is a simple exchange protocol which allows the integration of IP devices into ATM networks without the need for the IP devices to run PNNI at all. An IP device can be connected to the network via a PAR-enabled device which is configured as a Proxy-PAR server The IP device itself is configured as a Proxy-PAR client. In accordance with Proxy-PAR, the Proxy-PAR client can register details of the IP services it supports with the Proxy-PAR server. This information is then encapsulated in PAR PTSEs as previously described and flooded in the network in the usual PNNI way. The Proxy-PAR client can also request the Proxy-PAR server for information on other IP devices connected in the network for which PAR PTSEs have been received by the PAR-enabled device as previously described. In this way, IP information is communicated between IP devices without the devices participating in PNNI.\nThrough use of PAR and Proxy-PAR as described above, protocol devices, in particular IP devices, can learn about each other via this communication of protocol information over the PNNI network, avoiding the need for manual input in each device of the protocol information needed for configuration of the higher layer protocol topology. For example, IP routers at the edge of an ATM cloud can learn about each other, and manual configuration of the IP adjacencies can be avoided. Further, our copending European Patent Application No. 99115544.1, filed 6 Aug. 1999, discloses mechanisms for dynamic configuration of OSPF (Open Shortest Path First) interfaces in IP routers. Routers in mobile networks, for example, can dynamically configure OSPF interfaces with the OSPF area of other (fixed or mobile) network routers as the mobile network roams and makes new connections Whether or not OSPF interfaces are configured dynamically, PAR and Proxy-PAR allow routers to register their protocol information (e.g. IP address, ATM address, OSPF area) with their serving ATM switches which then flood the data throughout the network. Other routers can retrieve this IP information by querying their serving ATM switches. Routers can then exchange routing information to form neighbor relationships, or “peer”, in the usual way with other routers they learn about from the information received.\nThus, in general, when PAR is used for communication of protocol information between protocol devices over a PNNI network as described above, configuration of the higher layer protocol topology over the physical ATM network is based on the protocol information supplied by PAR-enabled switches to their client devices via Proxy-PAR. At present, a PAR-enabled switch responds to a Proxy-PAR request from its client device by supplying all protocol information of the type requested which has been received in PAR PTSEs from the network. The present invention is predicated on the realization that this is unnecessary and can compromise efficient configuration of the network topology for the protocol in question."}
{"text": "In transport refrigeration systems, such as refrigerated trucks, truck-trailers or refrigerated containers, during periods when the compressor is operating to compress the refrigerant in the system, substantial pressures can build up on the discharge (i.e. high pressure side) of the system. The vapor compression circuit is therefore designed to safely contain these pressures. It is recognized, however, that situations may arise where the pressures will tend to exceed what is considered a safe level. Accordingly, it is necessary to provide design features which will relieve these pressures before they become excessive.\nIn accordance with one established protocol, three levels of safety are provided on the high pressure side of the refrigeration system. The three levels are applied sequentially and in a prioritized fashion as follows. The first level is implemented in software and is based on pressure transducer readings. That is, when a predetermined pressure level is sensed, action is taken to limit the refrigerant flow, shut off the compressor or the system, or temporarily shut off the system and restart it after the pressure drops within a tolerance range.\nA second level is implemented by way of a mechanical pressure switch which responds to sensed pressures to shut the system off or temporarily shutting the system off and restart it after a period of time.\nA third level is implemented by way of a mechanical relief device which responsively opens to at least partially allow the refrigerant to be released to the atmosphere in the event that prescribed pressure levels are exceeded.\nRecently concerns have arisen about the environmental effects of the release of commonly used refrigerants to the atmosphere by way of leakage and the like. One approach to addressing this problem is the use of a more benign refrigerant, CO2, in place of the traditional refrigerants such as Freon. With such a refrigerant, however, it is necessary to operate at substantially higher pressures, and therefore compressors have been designed specifically for the compression of CO2. With these higher pressures in the circuit, it is even more important to continuously sense these pressures and when they become excessive, provide relief in a safe manner. For that purpose, the three level protocol as described above has been found to be satisfactory to control the operating pressures on the high pressure side during operation of the system.\nWith the use of CO2 as the refrigerant, the applicants have recognized that, in addition to the occurrence of excessive pressures during periods of operation, system pressures may also become excessive during periods of shipping and storing. That is, when a charged system at rest is exposed to excessive ambient temperature conditions, such as may occur in a warehouse in the summertime or when a system is exposed to direct sunlight at midday, the pressures are likely to rise to undesirable levels. Under these conditions, the three level safety protocol will be useful in relieving pressure on the high pressure side, but, unlike systems with conventional refrigerants in these situations, a CO2 system will be susceptible to excessive pressure conditions on the low pressure side as well.\nWhat is needed is therefore a method and apparatus for relieving pressures on the low pressure side of a CO2 system which is exposed to high ambient temperature conditions during shut down."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a projection apparatus, and particularly, to a projection apparatus having a plurality of light sources.\n2. Description of Related Art\nAs science and technology develop, demand of projection apparatus for higher brightness are increasing, which can not be satisfied by that having a single light source. Accordingly, projection apparatus employing multiple light sources have been proposed.\nHowever, projection apparatus employing multiple light sources correspondingly consumes more electricity, which also generates excessive heat. Therefore, the heat dissipation system, especially for the multiple light sources, becomes a critical concern in designing the projection apparatus. Light sources generate a lot of heat, thus heat dissipation systems for light sources are very important for controlling the temperature thereof within a range prescribed by the manufacturers in order to maintain the characteristics and lifetime of the light sources.\nIn a conventional projection apparatus employing multiple light sources, a typical heat dissipation system used for conventional projection apparatus employing a single light source is used, which comprises a blower, a wind duct and an axial fan disposed adjacent to the single light source, and the wind duct is capable of guiding a cooling airflow generated by the blower to a burner of the light source for cooling down the temperature thereof. The axial fan is capable of generating another cooling airflow for cooling a reflecting cover of the light source. However, when such a heat dissipation system is employed in a projection apparatus using multiple light sources, corresponding numbers of fans, wind ducts and axial fans are required, which increases not only the cost but also the volume thereof. Further, using more blowers not only consumes more electricity, but also makes the system circuit thereof more complex. Also, using more blowers even limits the size of the projection apparatus, so that the bulkiness thereof cannot be minimized."}
{"text": "With the advent of connection-oriented forwarding technologies such as Provider Backbone Transport (PBT) and Provider Backbone Bridge (PBB), Ethernet is rapidly becoming a dominant broadband technology, particularly in metro networks and wide-area networks. With PBT, service providers are able to establish point-to-point and point-to-multipoint Ethernet tunnels and to specify paths that service traffic will take through their Ethernet networks. With PBB, service providers are able to separate a communications network into customer domains and service provider domains. The separation is achieved by encapsulating the customer packets within a backbone (i.e., service provider) MAC (Media Access Control) header. Network elements in the service provider domain forward packets based on the service provider MAC header while the customer header remains invisible except at the edge of the service provider domain.\nAs Ethernet services proliferate, service providers require a robust set of operation, administration, and maintenance (OAM) tools to manage their Ethernet networks and to adapt the Ethernet technology to a carrier-grade service environment. To this end, the IEEE (Institute of Electrical and Electronics Engineers) organization has formalized a standards document for connection fault management in Ethernet networks, referred to as IEEE 802.1ag (also known as Connectivity Fault Management or CFM). The ITU-T Recommendation Y.1731 also defines OAM functions and mechanisms for Ethernet-based networks much like the 802.1ag standard. In general, such standards specify managed objects, protocols, and procedures for, among other things, detecting and diagnosing connectivity faults in end-to-end Ethernet networks. Defined CFM mechanisms for fault detection include continuity check, linktrace (traceroute), loopback (ping), and alarm indication at different levels or domains (e.g., customer level, service provider level, and operator level).\nThe IEEE 802.1ag standard also defines various CFM entities and concepts, including maintenance domains, maintenance associations, and maintenance association end points. According to IEEE 802.1ag, a maintenance domain (MD) is “the network or the part of the network for which faults in connectivity can be managed”, a maintenance association end point (MEP) is “an actively managed CFM entity” that “can generate and receive CFM PDUs” (protocol data units or frames), a maintenance association (MA) is “a set of MEPs, each configured with the same MAID (maintenance association identifier) and MD Level, established to verify the integrity of a single service instance”, and a maintenance entity (ME) is “a point-to-point relationship between two MEPs within a single MA”. Additional details regarding such CFM entities are available in the IEEE 802.1ag/D8.1 draft standard, the entirety of which is incorporated by reference herein.\nIn metro Ethernet applications, connectivity across tunnels (also called connections) between MEPs is verified continuously through continuity check (CC) messages. A network element transmits such CC messages periodically at a variable interval, which can occur as often as once every 3 milliseconds. Typically, the generating and processing of such CC messages occurs centrally, that is, by a general-purpose central processing unit on a processor card in the network element. The line cards extract the frames of the CC messages from the data path and send them to the processor card. In effect, this frame extraction and forwarding concentrates the CC messages from all line cards at this central point.\nBecause many connections (e.g., PBB/PBT tunnels) can terminate on a given physical interface on the network element, the central processor can become overwhelmed by the real-time processing requirements for generating and checking these CC messages. For example, a network element that supports 640 G of service traffic and has a scaling requirement of 1000 MEs per 10 G lane can thus have 64000 MEs to manage, with the corresponding CC messages converging on the single central processor. With a minimum interval for a CC message being 3.1 ms, the central processor can conceivably need to generate a CC message every 48 ns. Even the fastest of today's CPUs would not measure up to the task. Consequently, the CPU would eventually lag behind with CC message generation and checking, thus eventually leading to false indicators of lost connectivity. Alternatively, multiple general-purpose CPUs can be used in parallel, but this configuration can be impractical with respect to area, power consumption, and cost."}
{"text": "One-time use of hollow hypodermic syringe needles has become the norm in most health facilities because they cannot be reliably sterilized for reuse. The safe disposal of used needles has however, created a problem of its own; the cost to society is enormous and the end results have been less than satisfactory. Bulk incineration of accumulated syringes and needles at the point of use is an obvious but impractical solution because it cannot be accomplished easily, effectively or even safely. And once started toward the dubious but expensive channels of hazardous waste disposal prior to incineration or burial at a central facility the destiny of the syringes with their needles intact or of the needle assemblies per se seems to be all but uncontrollable.\nMechanical shearing devices have long been used at the point of use to sever used needles from their supports. While shearing goes far to prevent reuse of the needle particularly if it is sheared in sections, it is far from a sanitary solution. The severed parts remain sharp as well as contaminated and the shearing operation itself is believed to spread vaporized contamination as part of the shearing impact.\nIt has been known from the beginning of the electrical age that metallic electrical conductors having a low ratio of volume to surface area, of which a hollow needle is a classic example, will be incinerated when large electric currents are passed through them. A typical modern day device of this type is shown in U.S. Pat. No. 4,628,169 in which portions of the hypodermic needle are successively heated in relatively short sections from the tip inward to achieve incineration temperatures between fixed, spaced apart electrodes while the needle is still attached to the syringe body. No known prior art, point-of-use devices of this type, however, eliminate sharp needle stubs close to the syringe body and none is able to generate heat for a sufficient time to sterilize needle parts much beyond the points of electrode contact. Predetermined or optimized spacing of the fixed electrodes is not a solution to the problem because widely spaced electrodes, in which the spacing approximates the length of the exposed shaft of the needle, leave relatively long sharp stubs of needles attached to the syringe body, for reasons explained below. In addition, the time interval during which current can flow is minimal, with the circuit often being burned open before there is time for heat to be conducted to sections inward of the electrode coupling points to the needle shaft. Close electrode spacing, on the other hand, also restricts the time for the heat to flow along the needle to parts outside of the electrodes, such as those within the syringe body and, in addition, subjects the electrodes to becoming permanently short circuited by fused needle parts to shut down the entire process. And the most minimal practical electrode spacing nevertheless leaves a critical sharp needle stub requiring careful handling to overcome the latent contamination risks."}
{"text": "In the prior art, a folded antenna 50, for example, shown in FIG. 10 has been introduced as an antenna suitable for transmitting and receiving signals of two or more frequency bands having approximately two times relationship (for example, cf. PATENT DOCUMENT 1). In the folded antenna 50, it is known that resonance at two or more frequency bands having approximately two times relationship can be achieved, since resonance frequencies can be adjusted by arranging intervals between adjacent segments 51, or the like. The antenna of this kind is formed on a surface of a cylindrical base made of, for example, dielectric materials or the like, mounted on a top part of a casing as an antenna for a portable telephone or the like, and used by setting the antenna element 51 located through a certain distance h to a ground face (a ground plate) 52 which is the casing of the portable telephone or the like.\nPATENT DOCUMENT 1: Japanese Patent Application Laid-Open No. HEI10-13135"}
{"text": "The present disclosure relates generally to uninterruptible power supplies (UPSs), and particularly to synchronization control systems for control thereof.\nUPSs are employed in a variety of applications where a constant source of power is desired at a load. A typical UPS system involves an inverter feed path, also generally referred to as the inverter, that is operably connectable in parallel with a bypass feed path, also generally referred to as the mains. The mains may be connected to a utility, but may also receive power from some other supply not connected to a utility electrical grid. The inverter may receive power from the same source as the mains, but may also receive power from some other supply.\nThere are several types of UPSs depending on their operation mode. Double conversion UPSs offer the maximal protection level as the load is always fed by the inverter. On the other hand, with line-interactive UPSs, the load is fed by the mains and the inverter is used to correct the shape of the load voltage.\nThere are also several possible UPS configurations to supply a critical load, such as Redundant Parallel Architecture (RPA), Dual Independent Configuration, Load Bus Synchronization, and Power Tie, for example. With the RPA concept, (N+M) UPSs are paralleled to supply a load that can be fed by N UPSs only. This way, a redundancy of M units is achieved. More and more, and for high availability, Dual Independent Configurations are requested by customers. This requires the synchronization of two independent UPS groups and the use of an Intelligent Static Switch (ISS) that automatically switches the critical load from one source to the other. Another concept is Load Bus Synchronization where two independent UPS groups can be temporarily synchronized in order to move the critical load from one side to the other for maintenance purposes. An extension of the Load Bus Synchronization concept is the Power Tie concept, where the two independent UPS groups are permanently synchronized and their load shared as if they were a unique UPS group in a RPA configuration. Finally, and with consideration to the bypass configuration, it is desirable to be able to choose between two different options, a centralized bypass or decentralized bypass.\nAccordingly, there is a need in the art for a control system and apparatus that allows multiple configurations of UPSs in critical power management systems."}
{"text": "The present invention pertains to a configurable system memory map. More particularly, the present invention pertains to the modification of a memory map so as to improve performance in a computer system or the like.\nIn a computer system or the like, a memory map may be created which reflects the addressability of devices coupled together in a system including a processor. For example, a computer system or the like may include a processor that accesses a variety of addressable components such as the following: a buffer memory, an internal and/or external Read-Only-Memory (ROM), one or more peripherals, internal or external static random access memory (SRAM), etc. In some cases both memory devices and input/output devices are accessed from the same memory map. In other words each address identifies either a location in an internal or external device, or identifies a different type of device such as a peripheral device (or a storage area in the peripheral device).\nThe memory map is especially important during reset and interrupt processing. As is known in the art, when a processor is reset, the processor will look for instructional code to execute at a certain memory location. For example, with a 32-bit address space, the processor may first access address 00000000 (hex) for the first instruction code to be executed so as to set up the processor to execute additional code (one skilled in the art may refer to this initial code as a xe2x80x9cboot-upxe2x80x9d or initialization sequence of code). Because a reset condition typically occurs after a power-down/power-up sequence, it is necessary that the first address accessed after reset be in non-volatile memory.\nMany processors also receive interrupt inputs which will cause the processor to stop execution of the code sequence currently being executed and to jump to another, preset address. In many processors, the preset address is a relatively low address like the reset address (e.g., the example above, 00000000(hex)). For example, an interrupt address could be 00000100 (hex). Then, the processor can execute code at this address. The processor may continue to execute instruction codes at sequential address locations, or the instruction code stored at the interrupt address may cause the processor to jump to another address for instruction codes to be executed to handle the interrupt. When execution of the code to handle the interrupt is completed, the processor returns to the execution of code that was previously interrupted.\nAs stated above, the reset address and the interrupt addresses (sometimes individually referred to as vectors) are typically located relatively close together at the lower or upper part of the address space. Since the reset vector must point to valid instruction code, it usually points to a non-volatile memory device. It is expected, then, that the interrupt vectors will also point to the same non-volatile memory device. Furthermore, the non-volatile memory may be an external device because on-chip, rewritable non-volatile memory may not be available in the IC technology used and the code may not be available at the time of fabrication of the chip to use internal ROM. A problem that is seen with this kind of configuration is that code execution on interrupts may occur with external, non-volatile memory that provides a slower access time when compared to other types of memory such as an internal SRAM device. In applications where interrupts must be handled at high speed, or must be dynamically changed during normal system operation, continuous access to the non-volatile memory space (e.g., an external ROM) during interrupt processing may not be acceptable. Though the first instruction(s) from the ROM may cause a jump to a different memory address (e.g., a memory address that points to a faster memory device), these initial accesses can waste processing time and negatively impact performance.\nIn view of the above, there is a need for an improved method and apparatus for a computer system using a memory map.\nAccording to an embodiment of the present invention, a configurable memory map is provided. During a first condition, such as during a reset, the memory map is in a first configuration so that access to a first portion of the memory space (e.g., for a reset vector) results in an access to a first memory device. During a second condition (e.g., during normal processing), the memory map is in a second configuration (e.g., through the operation of a bus controller) so that accesses to the same portion of the memory space result in accesses to a second memory device (e.g., during interrupt processing). Using such a configurable memory map can result in improved performance for the computer system. For example, accesses by a processor during a reset operation will result in accesses to non-volatile memory (e.g., a ROM) while accesses by the processor during interrupt processing will result in accesses to faster memory devices (e.g., an internal SRAM) even though the reset and interrupt vectors are in the same portion of the memory space. In addition, access to the first memory device and access to the second memory device are available in other fixed portion of the memory devices to facilitate initialization of the volatile memory device (e.g., internal SRAM) after reset."}
{"text": "High speed uplink packet access (HSUPA) is a technology introduced in the 3rd Generation Partnership Project (3GPP) Release 6 (R6). This technology is optimization and evolution directing to packet services in the uplink direction (direction from a mobile terminal to a radio access network). Using self-adaptation encoding, physical-layer hybrid retransmission, Node B-based fast scheduling, and 2-ms transmission time interval (TTI)-based short-frame transmission, HSUPA enhances the highest data transmission rate and cell throughput, and reduces the delay.\nHSUPA further reduces the transmission delay by using the 2-ms TTI transmission mode. However, 2-ms TTI-based short-frame transmission mode is not applicable to the user equipment (UE) that works at the edge of a cell. Therefore, the coverage of the UE is narrowed down. Especially in the case of a period of time after the UE's total transmit power reaches the maximum transmit power, the 2-ms TTI needs to be switched to 10-ms TTI to enhance the uplink coverage. The detailed process is as follows:\n(1) The network configures 6d event for the UE by a measuring control message.\n(2) When the UE transmit power reaches the maximum allowable transmit power and stays for a period of time, the UE reports a measurement report to the network.\n(3) Upon receiving the measurement report, a radio network controller (RNC) judges whether the UE uplink transmit power is limited, and sends a reconfiguration message (may include a radio bearer reconfiguration message, transmission channel reconfiguration message) to the UE to reconfigure the TTI length.\n(4) Upon receiving the reconfiguration message, the UE performs the reconfiguration according to the reconfiguration message at a time specified by the network to implement switching from 2-ms TTI to 10-ms TTI.\nAs described above, in the prior art, the UE reports 6d event and the RNC judges, according to 6d event, whether the UE transmit power is limited, and initiates the reconfiguration message to configure the UE for TTI switching. Therefore, the delay is long when the UE switches to 10-ms TTI transmission mode."}
{"text": "This invention relates to a brake operating device adapted to be mounted to a bicycle handle bar, and more particularly to a brake operating device adapted to be mounted to a bicycle drop type handle bar comprising a straight rod section and substantially U-shaped bent rod sections extending forwardly and downwardly from both ends of the straight rod section.\nGenerally, this kind of brake operating device is provided with a bracket member mounted to the bent rod section of the handle bar, the bracket member carrying a main lever substantially parallel to the bent rod section and an auxiliary lever substantially parallel to the straight rod section and bent rod section, so that a rider may operate the brake by hands gripping the handle bar at the position of either the straight rod section or the bent rod section.\nThe main lever supports one end of a control wire connected at the other end thereof to the bicycle brake and the auxiliary lever is made operable in association with the main lever. Such means for associating the levers are, as is well-known, disclosed in the British Pat. No. 778,809, and the U.S. Pat. No. 3,403,577 and No. 3,915,028. Namely, the first Patent discloses an auxiliary lever fixed to one side of a main lever through a screw bolt, the second discloses an auxiliary lever supported to an axial end of a pivot swingably supporting a main lever, the auxiliary lever extending for insertion at its head portion between the foremost end face of the main lever and a bracket member and being operable in association with the main lever, and the third discloses an auxiliary lever fixed at its head to a main lever through a single pin.\nHowever, in the first case of using the screw bolt, there is a defect such that the junction of both levers becomes weak. Also the main lever is subjected to the eccentric force of the brake operation of the auxiliary lever to thereby cause torsion on the pivot of the main lever, resulting in a brake action which is not smooth. In the second case, the defect is such that the main lever changes its position in the direction of the swinging motion whether or not the auxiliary lever is mounted to the bracket member, thereby considerably changing the distance between the handle bar and the grip of the main lever, and further the auxiliary lever is likely to be loose and unstable due to its one-point-support. In the third case, the problem is that the auxiliary lever is not easily removable and also unstable due to its one-point-support.\nIn order to overcome the aforesaid porblems, this invention has been designed. A main object of the invention is to provide a brake operating device capable of easily mounting or dismounting an auxiliary lever and of holding it in a stable manner. Another object of the invention is to provide a brake operating device in which the auxiliary lever is mounted at the rider's choice and which is available by the main lever without the auxiliary lever, and notwithstanding, the mounting of the auxiliary lever, the main lever position is kept constant in the swinging motion direction.\nThe brake operating device of the invention is so constructed that a hollow shaft is provided at the foremost end of the main lever supported swingably to a bracket member. The hollow shaft carries a retaining member holding one terminal of a control wire, the other terminal of which is connected to the bicycle brake. The auxiliary lever is mounted to the bracket member swingably around a pivot coaxial with a pivot through which the main lever swings and is provided at the head and with an engaging shaft parallel to the common axis of the pivots and insertable into the hollow shaft. The engaging shaft engages with the hollow shaft to thereby operate the auxiliary lever in association with the main lever.\nThe auxiliary lever is pivoted at its pivoting portion to the bracket member through the second pivot coaxial with the first pivot carrying the main lever, and is supported at the head insertably into the hollow shaft through the engaging shaft, thus being stable through a two-point-support.\nFurthermore, the auxiliary lever, which is in association with the main lever through insertion of the engaging shaft into the hollow shaft, is easily mountable or removable to or from the bracket member, and also usable at the rider's choice. Furthermore, the auxiliary lever, when in use, is inserted at its head not between the bracket member and the front end face of the main lever as in the aforesaid prior art, but into the hollow shaft for operation in association with the main lever, resulting in the position of the main lever in the swinging motion direction being kept constant regardless of the existence of the auxiliary lever.\nThese and other objects and novel features of the invention will be apparent from the following description of an embodiment of the invention in accordance with the accompanying drawings."}
{"text": "All publications mentioned throughout this application are fully incorporated herein by reference, including all references cited therein.\nOne of the major physiological roles of the endothelium is to preserve the integrity of the vasculature by its permeability barrier properties, and to provide a non-thrombogenic surface. Therefore, the endothelial cell surface serves as a prime regulatory site of coagulation responses. Under normal circumstances, an injury to vascular endothelial cells lining a blood vessel triggers a haemostatic response through a sequence of events commonly referred to as the “coagulation cascade”. The cascade culminates in the conversion of soluble fibrinogen to insoluble fibrin which, together with platelets, forms a localized clot or thrombus which prevents extravasation of blood components. Wound healing can then occur followed by clot dissolution and restoration of blood vessel integrity and flow.\nThe events which occur between injury and clot formation are carefully regulated and linked series of reactions, involving number of plasma coagulation proteins in inactive proenzyme forms and cofactors circulate in the blood. Active enzyme complexes are assembled at an injury site and are sequentially activated to serine proteases, with each successive serine protease catalyzing the subsequent proenzyme to protease activation. This enzymatic cascade results in each step magnifying the effect of the succeeding step. In brief, the blood coagulation cascade is usually initiated as soon as the cell surface glycoprotein tissue factor (TF) comes into contact with circulating activated factor VII (VIIa), resulting in the formation of TF-FVIIa complex. The TF-VIIa complex activates X to factor Xa. Subsequently, factor Xa catalyzes the conversion of prothrombin to thrombin, thereby leading to fibrin formation, platelet activation, and, ultimately, generation of a thrombus. Tissue factor pathway inhibitor (TFPI) is a potent direct inhibitor of factor Xa, and in a factor Xa dependent fashion, produces inhibition of the factor VIIa-TF complex.\nWhile efficient clotting limits the loss of blood at an injury site, inappropriate formation of thrombi in veins or arteries is a common cause of disability and death. Abnormal clotting activity can result in and/or from pathologies or treatments such as myocardial infarction, unstable angina, atrial fibrillation, stroke, renal damage, percutaneous translumenal coronary angioplasty, disseminated intravascular coagulation, sepsis, gestational vascular complications, pulmonary embolism and deep vein thrombosis. The formation of clots on foreign surfaces of artificial organs, shunts and prostheses such as artificial heart valves is also problematic.\nApproved anticoagulant agents currently used in treatment of these pathologies and other thrombotic and embolic disorders include the sulfated heteropolysaccharides heparin and low molecular weight heparin (LMWH). These agents are administered parenteral and can cause rapid and complete inhibition of clotting.\nHeparanase is an endo-β-D-glucuronidase capable of cleaving heparan sulfate (HS) side chains at a limited number of sites, yielding HS fragments of still appreciable size (˜5-7 kDa) [Freeman, C. and Parish, C. R. Biochem. J. 330:1341-1350 (1998); Pikas, D. S. et al. J. Biol. Chem. 273:18770-18777 (1998)]. Expression of heparanase is restricted primarily to the placenta, platelets, keratinocytes, and activated cells of the immune system, with little or no expression in connective tissue cells and most normal epithelia [Parish, C. R. et al. Biochem. Biophys. Acta 1471:M99-108 (2001); Vlodaysky, I. and Friedmann, Y. J. Clin. Invest. 108:341-347 (2001)]. During embryogenesis, the enzyme is preferentially expressed in cells of the developing vascular and nervous systems. Heparanase activity was implicated in tumor growth, neovascularization, inflammation and autoimmunity [Parish (2001) ibid.; Vlodaysky (2001) ibid.; Dempsey, L. A. et al. Trends Biochem. Sci. 25:349-351 (2000)]. A single human heparanase cDNA sequence was independently reported by several groups [Hulett, M. D. et al. Nat. Med. 5:803-809 (1999); Kussie, P. H. et al. Biochem. Biophys. Res. Commun. 261:183-187 (1999); Toyoshima, M. and Nakajima, M. J. Biol. Chem. 274:24153-24160 (1999); Vlodaysky, I. et al. Nat. Med. 5:793-802 (1999)]. Thus, unlike the large number of proteases that can degrade polypeptides in the extra-cellular matrix (ECM), one major heparanase appears to be used by cells to degrade the HS side chains of HS proteolycans.\nHeparanase is synthesized as a latent 65 kDa precursor whose activation involves proteolytic cleavage at two potential sites located at the N-terminal region of the molecule (Glu109-Ser110 and Gln157-lys158), resulting in the formation of two protein subunits, 8 and 50 kDa polypeptides, that heterodimerize and form the active heparanase enzyme [McKenzie, E. et al. Biochemical J. 373: 423-35 (2003); Levy-Adam F. et al. Biochem. Biophy. Res. Commun. 308: 885-91 (2003)]. One of the prime physiological sources for heparanase are platelets [Hulett (1999) ibid.; Freeman and Parish (1998) ibid.]. The 50 and 8 kDa heparanase polypeptides were biochemically purified from platelets, which also contain significant amounts of the 65 kDa proenzyme [Hulett (1999) ibid; Freeman (1998) ibid.]. Moreover, the heparanase gene was previously cloned from human platelets [Hulett (1999) ibid.]. Heparanase released by activated platelets or platelet-derived microparticles, is biologically active, stimulates angiogenesis and modulates endothelial cell activities [Brill, A. et al. Cardiovasc. Res. 63:226-235 (2004); Myler, H. A. and West, J. L. J. Biochem. 131:913-922 (2002)].\nIt has been previously shown by applying heparanase which lacks enzymatic activity that heparanase exerts also non-enzymatic activities, independent of its involvement in ECM degradation and alterations in the extracellular microenvironment associated with angiogenesis, cell survival, and migration [Goldshmidt, O. Faseb. J. 17:1015-1025 (2003); Gingis-Velitski, S. et al. J. Biol. Chem. 279:23536-23541 (2004); Zetser, A. et al. Cancer Res. 66:1455-1463 (2006)].\nThe inventors have recently demonstrated that heparanase is also involved in the hemostatic system. Heparanase was shown to up-regulate TF expression [Nadir, Y. J. et al. Thromb. Haemost. 4:2443-2451 (2006)] and interact with TFPI on the cell surface, leading to dissociation of TFPI from the cell membrane and increased cell surface coagulation activity [Nadir, Y. et al. Thromb. Haemost. 99:133-141 (2008)].\nWhen the present inventors further examined the role of heparanase in coagulation modulation, they surprisingly and unexpectedly found, as shown by the present invention, that heparanase is directly involved in the activation of the coagulation system by enhancing factor Xa production in the presence of TF/VIIa complex.\nMore particularly, addition of heparanase to TF and factor VIIa resulted in three to four fold increase in activation of the coagulation cascade as depicted by increased factor Xa and thrombin production. When heparanase was added to pooled normal plasma a seven to eight fold increase in Xa level was observed. Subsequently, the inventors looked for clinical data supporting the existence of this innovative role of heparanase. Plasma samples of thirty five patients with acute leukemia at presentation, and twenty healthy controls were studied for heparanase and factor Xa levels using ELISA and chromogenic assay, respectively. A strong positive correlation was found between plasma heparanase and Xa levels (r=0.735, p=0.001). In addition, unfractionated-heparin and an anti-Xa derivative abolished the effect of heparanase, while, TFPI and TFPI-2 only attenuated the procoagulant effect. Finally, co-immunoprecipitation and far-western analyses demonstrated that heparanase directly interacts with TF. Altogether, the results presented herein support the notion that heparanase is a modulator of blood homeostasis, and suggest a novel mechanism by which heparanase is involved in direct activations of the coagulation cascade by direct interaction with TF.\nFor the purpose of better assessing the procoagulant role of heparanase in various medical conditions, the inventors have developed an original assay to evaluate heparanase procoagulant activity in biological samples. Heparanase procoagulant activity was preliminary studied using purified proteins of heparanase, TF, factor VIIa and factor X. The methodology was verified in fifty five plasma samples and compared to heparanase and tissue factor pathway inhibitor (TFPI) levels by ELISA and factor Xa, thrombin levels and antithrombin activity by chromogenic substrates. Thirty five samples were of third-trimester pregnant women (weeks 36-41) who were in labor or came for appointed elective cesarean section and 20 control samples were of non-pregnant healthy women. The heparanase procoagulant activity assay was shown to differentiate heparanase procoagulant effect from TF activity, in purified proteins model. Heparanase procoagulant activity was significantly higher in the plasma of pregnant women compared to non-pregnant (p<0.005). Heparanase relative contribution to the TF heparanase complex activity was significantly higher in the plasma of pregnant women compared to non-pregnant (29% increase, p<0.0001). Differences in heparanase procoagulant activity were more prominent than changes in heparanase levels by ELISA, TF activity, factor Xa, thrombin and free TFPI levels. The results presented therein demonstrate that heparanase procoagulant activity can be determined by the assay and the methodology of the present invention. In this particular case, the assay of the invention revealed a significant contribution of heparanase to the procoagulant state in late third-trimester pregnancy and at delivery.\nThe methodology and the assay developed in the present invention were then applied to the field of orthopedic hip and knee surgeries which are followed by a hypercoagulable state. Heparanase level and procoagulant activity in patients undergoing orthopedic surgery were assessed. The study group included fifty orthopedic patients. Thirty one patients underwent hip surgery and nineteen had knee operation. Fifteen individuals suffered from traumatic hip fractures and thirty five had osteoarthrosis of hip or knee joints. All patients received prophylactic dose of enoxaparin (a low molecular weight heparin) starting 6-8 hours post operation and lasting for five weeks. Plasma samples were drawn preoperatively and at 1 hour, 1 week and 1 month post operation. Samples were tested for heparanase levels by ELISA and TF/heparanase complex activity, TF activity, heparanase procoagulant activity (with the assay of the present invention), factor Xa and thrombin levels using chromogenic substrates. Heparanase levels were significantly higher 1 hour and 1 week post operatively compared to preoperative levels (p<0.05, p<0.005, respectively). The most dramatic changes were observed in heparanase procoagulant activity reaching a 2 fold increase 1 week postoperatively and 1.7 fold increase 1 month after surgery (p<0.0001, p<0.0001, respectively). Levels of factor Xa and thrombin did not significantly change. In conclusion, the inventors demonstrated that heparanase is involved in coagulation activation of orthopedic surgery patients. Heparanase procoagulant activity is highest 1 week postoperatively and remains high 1 month after operation. In view of the above, considering extending prophylactic anticoagulant therapy or evaluating heparanase procoagulant activity may potentially prevent late thrombotic events.\nIt is therefore one object of the invention to provide a composition comprising a combination of heparanase or any fragments, derivative or peptides thereof and TF or any fragments, derivatives or peptides thereof, specifically for use in modulation of coagulation processes.\nIn yet another object, the invention provides a method for the activation of a pro-coagulation process in a subject in need thereof, using the combined heparanase-TF compositions of the invention. These methods are specifically applicable in the treatment, amelioration and prevention of coagulation-related pathologic conditions\nIn yet another object, the invention provides the use of the novel heparanase-TF interaction as a target in a screening method for identification of coagulation-modulating compounds. Such compounds modulate the heparanase-TF mediated coagulation activity in a subject suffering of a coagulation-related disorder.\nIn yet another object, the invention provides a method for determining heparanase procoagulant activity in a biological sample of a mammalian subject.\nIn yet another object, the invention provides a diagnostic method for the detection and/or monitoring of a coagulation-related pathologic disorder in a mammalian subject.\nIn yet another object, the invention provides a diagnostic kit for the detection and/or monitoring of a coagulation-related pathologic condition in a mammalian subject.\nThese and other objects of the invention will become apparent as the description proceeds."}
{"text": "Voltage-gated calcium (Ca2+) channels (VGCC) play an integral role in the regulation of membrane ion conductance, cellular excitability and neurotransmitter release. VGCC are composed of the pore-forming α1 subunit and auxiliary α2δ ppm and β subunits that modulate channel expression and function. Among the low voltage activated channels is the Cav3 channel subtype, which mediates T-type calcium currents that may be targeted for treatment of epilepsy, especially children absence epilepsy and chronic pain (Huguenard, 1998, Cribbs et al., 2000, Perez-Reyes et al., 2009, Perez-Reyes, 2010).\nThe T-type or “low voltage-activated” calcium channels are so named because they open for shorter duration (T=transient) than the L-type (L=long lasting) calcium channels. T-type channels are activated at relatively negative membrane potentials (˜−60 mV). In many types of neurons, Ca2+ influx through T channels triggers low-threshold Ca2+ spikes, which in turn elicit a burst of action potentials mediated by voltage-gated sodium (Na+) channels. Brief burst firing is thought to play an important role in the synchronized activity of the thalamus and neuronal pacemaker under physiological conditions, but it also underlies a wide range of thalamocortical dysrhythmias under pathological conditions such as neuropathic pain or seizures, T channels can be activated by mild depolarization of the cell membrane (Talley et al., 1999, Perez-Reyes, 2003, Perez-Reyes, 2010, Pexton et al., 2011, Todorovic and Jevtovic-Todorovic, 2011).\nMolecular cloning has revealed three distinct T channel proteins, designated Cav3.1, Cav3.2 and Cav3.3. The Cav3.1 and Cav3.3 channels are expressed predominantly, though not exclusively, in the CNS. In contrast, the Cav3.2 channel is not only present in the CNS, but also expressed in peripheral nerve cell bodies and nerve endings of afferent fibers (Huguenard, 1998, Cribbs et al., 2000, Perez-Reyes et al., 2009, Perez-Reyes, 2010). The Cav3.2 channel is highly expressed in dorsal root ganglion (DRG) neurons, whereas little Cav3.1 and virtually no Cav3.3 are expressed in the small diameter DRG neurons (Nelson et al., 1992). The Cav3.2 channels are also expressed at a lower level in several non-neuronal tissues, including heart, liver, kidney, and pituitary. Both diabetic neuropathy and chronic constriction injury models in rats lead to DRG neuron-specific upregulation of the Cav3.2 channel and the T current density. This pathological adaptation results in enhanced excitability of sensory neurons and causes hyperalgesia and allodynia (Jagodic et al., 2007, Jagodic et al., 2008, Latham et al., 2009, Messinger et al., 2009, Yue et al., 2013). Conversely, knockout or antisense knockdown of the Cav3.2 isoform produces analgesic effects (Messinger et al., 2009).\nT-type channel inhibitors have two known uses in the clinic. The anti-absence seizure effects of ethosuximide and lamotrigine are thought to be mediated by the inhibition of T channel activity in the thalamus (Gomora et al., 2001, Huguenard, 2002). However, both drugs are weak and not specific against the T channel (Xie et al., 1995, Zhang et al., 1996). The antihypertensive effect of mibefradil is conventionally attributed to its inhibition of the T channel. However, mibefradil has poor selectivity with about 3-10 times more potent inhibition of the T-type than of the L-type Ca2+ current or the voltage-gated Na+ current (Avdonin et al., 2000). Because there are no selective T channel blockers, it is unclear whether and to what extend the inhibition of T channel activity at therapeutically relevant concentrations contributes to the therapeutic usefulness of a wide range of drugs, such as analgesics, antiepileptics, neuroprotectants, antipsychotics, antidepressants, antiarrhythmics and antihypertensives.\nTargeting a T-channel, particularly the Cav3.2 isoform, would be highly useful in reduction of thermal hyperalgesia and mechanical allodynia under pathological conditions, for example diabetic neuropathy. Several efforts to discover potent and selective T-type Ca2+ channels have been described in the literature, as exemplified below.\n1,4-Substituted piperidines, for example, “compound 30” (3,5-dichloro-N-{[1-(3,3-dimethylbutyl)-3-fluoropiperidin-4-yl]methyl}benzamide) and “TTA-P2” (3,5-dichloro-N-((1-((2,2-dimethyltetrahydro-2H-pyran-4-yl)methyl)-4-fluoropiperidin-4-yl)methyl)benzamide) were synthesized by Merck and found to potently block the T-Type Cav3.2 channel [J. Med. Chem. 51, 3692, (2008); J. Med. Chem. 51, 6471, (2008); US 2010/0222387; US 2013/8501773]. TTA-A2 suppresses active wake, promotes slow-wave sleep (Kraus et al., 2010), and prevents weight gain in mice on a high-fat diet (Uebele et al., 2009).\nA scaffold hopping approach afforded ML218 (3,5-Dichloro-N-[[(1α,5α,6-exo,6α)-3-(3,3-dimethylbutyl)-3-azabicyclo[3.1.0]hex-6-yl]methyl]benzamide, CID 45115620) a selective T-Type Ca2+ inhibitor. ML218 possess acceptable in vivo rat PK and was efficacious in a preclinical Parkinson model. Thus, ML218 is a useful new biologic probe to study T-Type Ca2+ function in vitro and in vivo (Xie et al., 2010, Xiang et al., 2011).\nCertain lactam acetamides have been described by Abbott and others as Cav2.2 and Cav3.2 calcium channel blockers, and ABT-639 has been reported as a Cav3.2 calcium channel blocker for treatment of diabetic neuropathic pain through peripheral action, because ABT-639 is presumed to not penetrate the blood brain barrier (Jarvis et al., 2014).\nN-Piperidinyl acetamide derivatives as calcium channel blockers have been described by Zalicus Pharmaceuticals, Ltd. [U.S. Pat. No. 8,569,344 (2013); U.S. Pat. No. 8,377,968 (2013)]. A piperidine-based compound, Z944, inhibits Cav3 channels in a voltage-dependent manner and is able to attenuate thalamic burst firing and suppress absence seizures in rats (Tringham et al., 2012). Z944 has shown promising results in clinical Phase I studies of pain in humans (Lee, 2014).\nDespite the fact that many T-type Ca2+ channel inhibitors have been discovered and have advanced to different stages of development, no FDA-approved selective T-type channel inhibitory compounds are available for clinical applications."}
{"text": "In some operating systems (e.g., version 9 of Mac OS from Apple Computer, Inc.), application programs (e.g., a web browser, a word processing software program, a game software program, or others) can assume that there is at least one frame buffer available for drawing their graphical user interfaces. A frame buffer is typically under control of a graphics controller (display controller), which controls a display device according to the information stored on the frame buffer. For example, a typical frame buffer is on the video memory on a graphics card. It is understood that drawing graphical user interfaces on a frame buffer refers to defining display information (e.g., in terms of RGB color values for a color display), typically pixel by pixel, for the display of the graphical user interfaces. The display controller uses the contents of the frame buffer to generate control signals for refreshing the screen of a display device (e.g., a Cathode Ray Tube (CRT) monitor, or a Liquid Crystal Display (LCD) panel).\nSome computers may not have a graphics controller and a display device. During the start up time of such a computer that has no graphics controllers, the operating system allocates a virtual frame buffer, using a portion of the system memory of the computer, to provide a place for the application programs to draw their user interfaces. The virtual frame buffer is not associated with or controlled by a particular graphics controller. The virtual frame buffer allows the application programs to generate the display information about their graphical user interfaces properly, even when the computer does not have directly attached graphics controllers and display devices. The content of the virtual frame buffer is not used for the generation of control signals for a display device that is directly attached to this computer."}
{"text": "1. Field of the Invention\nA method and apparatus for mounting a body part on a vehicle body for painting. In particular, a reinforcement member is connected to an exterior skin of the body part to prevent sagging of the body part during manufacture and subsequent use of the vehicle.\n2. Description of Related Art\nA vehicle body is the shell to which body parts of the vehicle are attached. Portions of the vehicle body that are easily visible after the vehicle is fully assembled include, for example, the A-pillar, the vehicle roof, and the areas surrounding the hood, doors, and trunk. Body parts typically include exterior vehicle components that are bolted onto the vehicle body. These body parts, which are non-integral to the body may include, for example, plastic garnishes, bumpers, fenders, grilles, passage doors, and fuel doors.\nTraditional body parts may be made of a plastic derivative or of metal. Plastic body parts and metal body parts each have inherent advantages and disadvantages. For example, while the plastic body parts may be formed into many more shapes than the metal body parts, the plastic body parts are prone to bending when a force is applied to them. This is undesirable to the owner of the vehicle and may cause the vehicle owner to have negative feelings about the vehicle. Specifically, the owner may believe that the vehicle is made with defective or low quality components. Further, plastic body parts are more sensitive to elevated temperatures that can occur during the vehicle manufacturing process, as will discussed hereinafter. On the other hand, while metal body parts are sturdy, a metal body part typically weighs more than the equivalent plastic body part. As such, the fuel economy for vehicles with a significant number of metal body parts is reduced. Metal body parts also may rust or corrode, thereby detracting from the appearance of the vehicle. Accordingly, a vehicle made of both metal and plastic body parts is desirable for a variety of reasons.\nDuring the manufacture of a vehicle, vehicle bodies and body parts are coated during a series of steps as they are conveyed along a coating line. This coating may be traditional spray painting, or may be for example, an electro-static coating. These coating steps include applying a primer to prepare the respective surfaces of the body and body parts for painting, applying a base coat for adding color to the body and body parts, and applying a clear-coat that functions as a protectant for the body and body parts.\nIn between these coating steps, the vehicle body and body the parts are passed through bake ovens. In the bake ovens, heated air is directed toward the body and the body parts. The bake ovens ensure that the coatings applied to the body and the body parts are adequately dried.\nWhen the body and the body parts are made of different materials, there are usually separate bake ovens for baking the body and the body parts. The separate bake ovens are usually operated at dissimilar temperatures. Furthermore, the amount of time that the coated body and coated body parts are in the different bake ovens can be altered to affect the amount of heat that is exposed to the coated body and coated body parts. For example, the coated body may be baked by exposure to a high temperature for a short period of time. Alternatively, the coated body part may be baked by exposure to a low temperature for a long period of time.\nThe body parts and the body that are made of dissimilar materials are separately passed through the bake ovens because of differences in the heat tolerances of the materials that make up the body and body parts. For example, the vehicle body is made of steel, and the body parts are made of some derivative of plastic. The steel vehicle body may be exposed to a much higher temperatures in the bake ovens than the plastic body parts. Unfortunately, there are several unintended consequences of preparing the parts in this separate manner.\nSince the vehicle body and the body parts are separately coated and baked, there can be an appearance mismatch between the vehicle body and the body parts. Specifically, because of lot to lot variation in coating ingredients and exposure to different baking conditions, the vehicle body and the body parts can have different appearances due to color and finish differences. The body and body parts may have slightly different color hues after being coated, the amount of sheen of the clear coat may be different, and/or the respective surfaces may have different levels of smoothness. As the body and the body parts are usually assembled so that they border each other, any appearance difference is easily detected and thus, undesirable.\nFurthermore, additional labor is required when the vehicle body and the body parts are separately coated. The body parts must first be coated and then later assembled with the vehicle body. As such, to be coated, the body parts are attached to a jig to be fed through the coating line and then removed from the jig after completion of the coating. Then, the body parts are installed on the vehicle body.\nAccordingly, it is considered apparent why it is not desirable to separately coat the vehicle body and body parts. However, numerous difficulties are also encountered when the vehicle body and body parts are coated and baked as a single entity. For example, as mentioned hereinbefore, the vehicle bodies are usually made primarily of some form of steel and the body parts can be made of a material different than steel, such as, for example plastic. When the combined body and body parts are coated and baked together, difficulties arise during the baking process.\nThe body parts, which are preferably made of plastic, are not typically as tolerant of high temperatures as the body that is made of steel. As such, damage can occur to the body parts when the body and body parts are jointly baked at a higher temperature. Primarily, any body part that has a horizontal section is especially susceptible to sagging of the body part. Sagging is not only an appearance defect, but can also cause misalignment issues. For example, if the body part is located near the vehicle hood and headlight opening, the sagging body part can prevent opening/closure of the hood.\nTherefore, there exists a need in the art for a body part that is made from a plastic material and is tolerant of the increased temperatures of a bake oven that is traditionally used for metal vehicle bodies."}
{"text": "This invention relates to methods for disposing of high-water sludge materials, particularly to methods of electrically dewatering sewage sludge.\nA variety of methods have been developed for disposing of sewage. In particular, a variety of biodegrading schemes have been perfected. Most of these schemes produce, as a final product, a sludge which is refractory to further biodegradation. The sludge consists of inorganic matter, proteins, cellulosic material, hydrocarbons and the like. Disposal of this sludge is a major problem for large sewage treatment facilities that may generate up to 100 tons a day of sludge.\nOne of the problems with sludge is that although it may appear to be very viscous or nearly solid, it is about 80 percent by weight water. Therefore, it is virtually non-combustible.\nOne method for dewatering sludge is what is herein referred to as electro-dewatering. In this process and electric field is applied to a mass of sludge. This tends to free the water, presumably because the constituent molecules of the sludge move in response to the electric field.\nFremont in U.S. Pat. No. 4,671,874 discloses an electro-dewatering system using a chopped and time varied current. The conveyors shown are not electrified, but have electrodes behind them.\nRoslonoski in U.S. Pat. No. 4,655,932 adds an electrolyte to the sludge to increase efficiency.\nBastegen in U.S. Pat. No. 4,402,834 teaches a conveyor which dewaters by vacuum, centrifugal, static pressure and pressure means.\nSneek in U.S. Pat. No. 4,101,400 teaches a pair of electrode conveyors inside a pair of motive conveyors."}
{"text": "Concrete slabs on grade, and other concrete structures and objects, traditionally are made by building a form or a mold. The forms and molds are usually made from wood, plywood, metal and other structural members. Unhardened (plastic) concrete is placed into the space defined by opposed spaced form members. Once the concrete hardens sufficiently, although not completely, the forms are removed leaving a concrete wall or other concrete structure, structural member or concrete object exposed to ambient temperatures. For some applications, such as runways, roads and highways, slip forms are used to continuously place long lengths of concrete. For both conventional forms, molds and slip forms, the unprotected concrete slabs, structures or objects are exposed to the elements during the remainder of the concrete curing process after the forming process. The exposure of the concrete to the elements, especially temperature variations, often makes the curing of the concrete a slow process and the ultimate strength difficult to control or predict. To compensate for these losses, larger amounts of portland cement are used than otherwise would be necessary. Larger cement fractions, particularly for high surface area-to-volume ratio elements, such as these, can increase the likelihood of shrinkage, cracking, and curling.\nThe curing of plastic concrete requires two elements, water and heat, to fully hydrate the cementitious material. The curing of plastic concrete is an exothermic process. This heat is produced by the hydration of the portland cement, or other hydraulic cements, that make up the concrete. Initially, the hydration process produces a relatively large amount of heat. As the hydration process proceeds, less hydration occurs thereby resulting in the production of less heat. At the same time, moisture in the concrete is lost to the environment. If one monitors the temperature of concrete during the curing process, it produces a relatively large increase in temperature which then decreases over time. This chemical reaction is temperature dependent. That is, the hydration process, and consequently the strength gain, proceeds faster at higher temperature and slower at lower temperature. During conventional concrete curing, first, the heat is lost; then, moisture loss makes it difficult for the cementitious material to fully hydrate, and, therefore, the concrete may not achieve its maximum strength.\nConcrete in conventional concrete forms or molds is typically exposed to the elements. Conventional forms or molds provide little insulation to the concrete contained therein. Concrete placed in forms for slabs on grade, roads, runways, and the like, is usually placed on a sheet of polyethylene placed on the ground first. Therefore, the bottom of the concrete slab is in contact with the ground, which absorbs the heat of hydration from the concrete. The top of the concrete slabs are sometimes covered with plastic sheathing to prevent some of the moisture loss to the environment. This also allows the concrete to lose the heat of hydration from the top of the slab surface to the air. Although during winter an insulated blanket may be placed on top of the concrete slab, road or runway to keep the concrete from freezing, heat produced within the concrete form or mold due to the hydration process usually is lost to the ground or to the air through a conventional concrete form or mold relatively quickly. Thus, the temperature of the plastic concrete may initially rise 20 to 40° C. or more above ambient temperature due to the initial hydration process and then fall relatively quickly to ambient temperature, such as within 12 to 36 hours. This initial relatively large temperature drop may result is concrete shrinkage, curling and/or concrete cracking. The remainder of the curing process is then conducted at approximately ambient temperatures, because the relatively small amount of additional heat produced by the remaining hydration process is relatively quickly lost to the ground or to the air through the uninsulated concrete form or mold. The concrete is therefore subjected to the hourly or daily fluctuations of ambient temperature from hour-to-hour, from day-to-night and from day-to-day. Curing concrete under ambient temperature conditions is not as significant a problem during summer temperature. It is cool or cold weather conditions that cause the most significant trouble for properly curing concrete.\nFailure to cure concrete under ideal temperature and moisture conditions affects the ultimate strength and durability of the concrete. In colder weather, concrete work may even come to a halt since concrete will freeze, or not gain much strength at all, at relatively low temperatures. By definition (ACI 306), cold weather conditions exist when “ . . . for more than 3 consecutive days, the average daily temperature is less than 40 degrees Fahrenheit and the air temperature is not greater than 50 degrees Fahrenheit for more than one-half of any 24 hour period.” Therefore, in order for hydration to take place, the temperature of concrete must be above 40° F.; below 40° F., the hydration process slows and at some point may stop altogether. It is typically recommended that concrete by moisture cured for 28 days to fully hydrate the concrete. However, this is seldom possible to achieve in commercial practice.\nInsulated concrete form systems are known in the prior art and typically are made from a plurality of modular form members. U.S. Pat. Nos. 5,497,592; 5,809,725; 6,668,503; 6,898,912 and 7,124,547 (the disclosures of which are all incorporated herein by reference) are exemplary of prior art modular insulated concrete form systems. Full-height insulated concrete forms are also known in the prior art. U.S. Patent Application Publication No. 2011/0239566 and applicant's co-pending patent application Ser. No. 13/247,133 filed Sep. 28, 2011 (the disclosure of which are both incorporated herein by reference in their entirety) discloses a full-height vertical insulated concrete form. Insulated concrete forms reduce the heat transmission to and from the concrete within such forms. However, vertical insulated concrete forms are not useful for forming slabs on grade, such as runways, road and highways and previously have not been proposed to be used in such applications.\nConcrete insulating blankets are known in the art. Electrically heated insulating blankets are also known in the prior art, such as those disclosed in U.S. Pat. Nos. 7,183,524 and 7,230,213. Such concrete insulating blankets and electrically heated insulating blankets are known for use in northern climates for thawing frozen ground and/or preventing curing concrete from freezing. It is known that plastic concrete will not cure satisfactorily at temperature below 50° F. However, such electrically heated blankets are designed to provide a constant amount of heat to the plastic concrete and are used only for the purpose of preventing the concrete from freezing in cold weather.\nTherefore, it would be desirable to produce a concrete forming or molding system for slabs on grade, such as runways, road and highways. It would also be desirable to provide a concrete curing system can be used specifically to cure slabs on grade, such as runways, roads and highways. It would also be desirable to provide a concrete curing system that accelerates concrete maturity or equivalent age to achieve improved concrete strength, particularly early concrete strength. Furthermore, it would be desirable to provide a concrete curing system for slabs on grade, such as runways, roads and highways, so that the concrete cures more quickly, is less permeable, more flexible, stronger, more durable and less prone to cracking and curling. It would also be desirable to produce concrete slabs on grade, such as runways, roads and highways, that are more environmentally friendly."}
{"text": "Over the years electrostatic and laser color copy/printers have shown significant improvement in their ability to make copies or prints giving excellent color rendition and image quality. The new generation of copiers and printers are now able to produce prints having quality comparable to that of silver halide color systems. With the advent of this ability the industry has attempted but has failed to produce, through electrostatic processes, images that have the look and feel of silver halide prints.\nPlain paper is typically used in electrostatic printing applications which does not generally provide a high degree of resolution, especially when color is involved. In photographic applications resin coated papers are used to provide the necessary resolution and quality. However, use of resin coated paper, i.e. polyethylene resin, as a copy or printing media in electrostatic printing applications has been a problem. Typical fuser roll temperatures are between 125 to 225.degree. C. Due to the low Tg of the resin the polyethylene softens or melts when coming into contact with the toner fuser roller of the copier/printer. This softening or melting causes paper jams and image degradation.\nIn addition to the \"melt problem\", electrostatic imaging directly on a resin coated paper substrate has been a problem. This is due to the toner used in electrostatic processes which is generally incompatible with the resin coated paper. Thus transfer and adhesion of the toner particles to the resin surface is not satisfactory and compromises the image production.\nThus in electrostatic imaging applications to have the feel and look of a standard silver halide print plain paper can not be used and use of \"photographic type\" substrates such as polyethylene resin or similar coated substrates are inadequate and pose a problem of softening or melting on the fuser roller. Accordingly, there is a need in the art for a resin coated paper for use in electrostatic applications to overcome these problems.\nThe invention provides such a solution by coating a coating layer over the resin coated substrate i.e. polyethylene layer. This coating layer typically comprises a natural or synthetic film forming polymer that has a melting point above 140.degree. C., thus preventing the resin coated substrate from melting and sticking to the fuser roller.\nThis protective layer provides a receiving surface resulting in photorealistic quality prints or copies. Depending on the optional ingredients in the layer a high gloss or semi-matte finish can be created. The invention provides advantage over conventional copying processes as well as over conventional photographic developing processes by providing an environmentally friendly process of producing photorealistic quality prints or copies without using toxic chemicals.\nGenerally, resin coated papers and adhesive-like gelatin \"subbing\" layers used for receiving an image are known in photographic applications.\nU.S. Pat. Nos. 3,811,913 and 4,188,220 to Kasugai et al. are representative of a \"subbing\" layer and resin coated paper, respectively, for use in photographic film processing. The '913 patent discloses use of gelatin and other polymeric materials as \"subbing\" layers for photographic materials including polyethylene coated paper substrates. UV radiation is applied to the coated surface to improve the adhesive property of the polyethylene support to the subbing layer. The '220 patent discloses a polyolefin coated paper. A low molecular weight polyolefin resin is incorporated into a conventional polyolefin resin to provide a coating layer.\nU.S. Pat. No. 4,312,937 to Kasper et al. discloses a resin coated paper including a paper layer and first and second layers of polyolefin adhered to opposite sides of the paper layer. Carbon black is incorporated into the polyolefin layers to eliminate pin-holing.\nU.S. Pat. No. 4,547,445 to Asahina et al. discloses a photographic material (\"postcard\") capable of having a photograph on one side and a writing surface of the opposite side. A paper support is coated with a polyolefin resin on both surfaces. A photographic emulsion layer is coated on one surface of the support and the opposite surface is coated with a gelatin layer including an inorganic pigment to absorb inks.\nU.S. Pat. No. 5,055,320 to Miura et al. discloses a support sheet including a subbing layer and a photographic emulsion layer. U.S. Pat. No. 5,075,196 to Daems et al. discloses supports for halftone dot image production. Daems provides a process including a paper base support coated on at least one surface with a polyolefin layer. On the exterior of the polyolef in layer is a white pigmented binder layer comprising a hydrophilic colloid binding agent and white titanium dioxide pigment particles. The light sensitive layer is coated on top of this binder layer.\nU.S. Pat. No. 5,082,724 to Katsura et al. discloses photographic paper supports consisting of a base paper support coated on both sides with a polyolef in resin.\nU.S. Pat. No. 5,104,721 to Sun discloses an electrophotographic imaging media comprised of a polymeric coating on a film substrate (slide projections) to improve printing resolution. The polymeric coating contains at least one pigment and has a Tukon hardness of from about 0.5 to about 5.0 and a glass transition temperature of from about 5 to 45.degree. C.\nU.S. Pat. No. 5,141,599 to Jahn et al. discloses a receiving material for ink jet printing including a polyolefin coated base paper with an ink receiving layer applied on the top surface. This receiving layer includes a mixture of gelatin and starch. The receiving material is defined as a gloss surface for ink jet printing comprising a polyolef in coated base paper and an ink receiving layer. The ink receiving layer contains a mixture of gelatin and starch in a ratio of 1:1 to 10:1 with the starch of a specific grain size.\nU.S. Pat. No. 5,182,161 to Noda et al., discloses a \"support for photosensitive materials\" comprising a base paper formed from natural kraft pulp according to a specifically defined digestion and chlorine bleaching process and a resin layer formed on the base paper. A subbing layer comprising a hydrophilic polymer such as gelatin is formed on the resin layer.\nIn addition the following Japanese patents all relate to papers suitable as \"photographic supports\" comprising paper coated on at least one side with a polyolefin and with one polyolefin surface over-coated with a hardener-containing gelatin layer. Japanese Patent No. 50-19402 dated Jun. 20, 1973; Japanese Patent No. 50-61154 dated Sep. 28, 1973; Japanese Patent No. 50-66519 dated Oct. 9, 1974; Japanese Patent No. 60-64306 dated Apr. 12, 1985; Japanese Patent No. 61-84643 dated Apr. 30, 1986; Japanese Patent No. 1-137252 dated May 30, 1989; and Japanese Patent No. 2-849 dated Jan. 5, 1990.\nIt is seen that gelatin over-coated polyethylene coated paper substrates in photographic applications is known. The Japanese patent abstracts describe such papers suitable as photographic supports, however, unlike the invention the gelatin layer is used as an undercoat on which a photographic emulsion coating is applied. This latter emulsion coating described in the prior art is the layer in which the image is formed by processing in an aqueous developing solution. These adhesive-like gelatin \"subbing layers\" are also described in Kasugai '913; Miura and Noda and also include a photographic emulsion coating thereon. Gelatin layers coated over polyolefin resin coated paper supports are described in Asahina et al. and Daems et al. However the exterior gelatin layer in Asahina includes an inorganic pigment and in Daems includes a white pigmented binder. The gelatin binder layer in Daems is further coated with a \"light sensitive\" layer.\nKasugai '220 and Kasper are representative of polyolefin coated papers. However, in contrast to the present invention, the resin layer in Kasugai '220 includes a low MW polyolef in resin and in Kasper includes carbon black. Both Katsura et al, and Sun disclose in general only \"polyolefin coated\" substrates: Kasura defines specific pulp fibers used to produce the base paper which is then coated with polyolefin; and Sun discloses a polymeric (film) substrate coated with polyolefin. Finally, Jahn discloses an ink jet sheet including a receiving layer which is a mixture of gelatin and starch.\nThe prior art does not teach a coated paper stock for electrostatic imaging. The present invention is directed to the provision of such by providing a substrate coated on at least one surface with a resin layer and a coating layer over the resin layer comprised of one or more natural or synthetic film forming polymers and has a glass transition temperature above 100.degree. C. The particular combination of resin coated paper of the invention provides an electrostatic copy or printing medium that is heat resistant. The coating layer of the invention also provides an image receiving surface layer in which toner particles are transferred and adhered to the surface during electrostatic imaging processes to produce photographic quality prints or copies.\nIt would be appreciated that advantage over known applications would be obtained by providing a coated paper where the melting point of the resin coating remains greater than 140.degree. C., preferably greater than 200.degree. C., to avoid problems of melting and toner adhesion during electrostatic imaging applications. In addition, this \"outer\" image layer provides a \"hard\" surface, which in a preferred embodiment comprises a gelatin and a crosslinking component. This is in contrast to the gelatin layers shown in the prior art which are soft and used as \"adhesive\" subbing layers. The \"hardness\" property is desired in the invention since the gelatin layer itself is used as the toner receiving surface layer.\nAccordingly, it is a broad object of the invention to provide a coated paper stock for electrostatic imaging comprised of a resin coated substrate with a coating layer.\nA more specific object of the invention is to provide a coated paper stock for electrostatic imaging where the \"outermost layer\" provides heat protection, gloss control, image improvement, improved smoothness, and improved toner adhesion and transport within the electrostatic imaging apparatus.\nAnother more specific object of the invention is to provide a coated paper stock for electrostatic imaging including a substrate coated on at least one surface with a resin layer comprised of olefinic material and an outer most heat protective layer, where the heat protective layer is a pin-hole free, continuous film over the resin layer.\nAnother object of the invention is to provide a coated paper stock and related process for producing a 3-dimensional relief image in an electrostatic imaging apparatus.\nAnother object of the invention is to provide a method for manufacturing a coated paper stock for electrostatic imaging.\nA specific object of the invention is to provide a dry method for producing photorealistic quality prints or copies which is advantageous over conventional photographic developing processes by being environmentally friendly and not using toxic chemicals.\nA more specific object of the invention is to provide a method for using a coated paper stock in electrostatic imaging applications without the associated problems of melting, image degradation and toner incompatibility."}
{"text": "In many applications a power converter is required to provide a voltage within a predetermined range formed from a voltage source having a different voltage level. Some circuits are subject to uncertain and undesirable functioning and even irreparable damage if supplied power falls outside a certain range. More specifically, in some applications, a precise amount of power is required at known times. This is referred to as regulated power supply.\nIn order to control a power converter to deliver a precise amount of power as conditions require, some form of control of the power converter is required. This control can occur on the primary side of an isolation transformer or the secondary side. A closed loop feedback control system is a system that monitors some element in the circuit, such as the circuit output voltage, and its tendency to change, and regulates that element at a substantially constant value. Control on the secondary side of a power converter can use a monitored output voltage as feedback control, but requires the use of some communication from the secondary to the primary side of the isolation transformer to control the primary side switching element. Control on the primary side can readily control the primary side switching element, but requires some feedback mechanism from the secondary side to the primary side to convey the status of the monitored element. In some applications, an optical coupler circuit, or opto coupler, is used to transmit feedback signals while maintaining electrical isolation between the primary and secondary sides.\nFIG. 1 illustrates a conventional regulated switch mode power converter including an optical coupler circuit. The power converter 2 is configured as a traditional flyback type converter. The power converter 2 includes an isolation transformer 4 having a primary winding P1 and a secondary winding S1. The primary winding P1 is electrically coupled to an input voltage Vin and a driving circuit including a transistor 8, a resistor 12, and a controller 10. A capacitor 28 is coupled across the input Vin and coupled with the primary winding P1. Input voltage to the circuit may be unregulated DC voltage derived from an AC supply after rectification and filtering. The transistor 8 is a fast-switching device, such as a MOSFET, the switching of which is controlled by the fast dynamic controller 10 to maintain a desired output voltage Vout. The controller 10 is coupled to the gate of the transistor 8. As is well known, the DC/DC conversion from the primary winding P1 to the secondary winding S1 is determined by the duty cycle of a pulse width modulation (PWM) switching signal provided to the transistor 8. The secondary winding voltage is rectified and filtered using the diode 6 and the capacitor 22. A sensing circuit and a load 14 are coupled in parallel to the secondary winding S1 via the diode 6. The sensing circuit includes a resistor 16, a resistor 18, and a secondary controller 20. The secondary controller 20 senses the output voltage Vout across the load.\nIn this configuration, the power converter is controlled by driving circuitry on the primary side, and the load coupled to the output is isolated from the control. As such, a monitored output voltage used for voltage regulation is required as feedback from the secondary side to the control on the primary side. The power converter 2 has a voltage regulating circuit that includes the secondary controller 20 and an optical coupler circuit. The optical coupler circuit includes two galvanically isolated components, an optical diode 24 coupled to the secondary controller 20 and an optical transistor 26 coupled to the controller 10. The optical diode 24 provides optical communication with the optical transistor 26 across the isolation barrier formed by the transformer 4. The optical coupler circuit in cooperation with the secondary controller 20 provides feedback to the controller 10. The controller 10 accordingly adjusts the duty cycle of the transistor 8 to compensate for any variances in an output voltage Vout.\nHowever, the use of an optical coupler circuit in and of itself presents issues. Firstly, the optical coupler circuit adds extra cost. In some applications, the optical coupler circuit can add more cost to the power converter than the isolation transformer. The optical coupler circuit also adds to manufacturing and testing costs. Furthermore, the performance of the optical coupler circuit degrades over time and therefore introduces another potential point of failure in the overall power converter. Also, characteristics of the optical coupler circuit must be accounted for in the overall circuit design. For example, the optical diode component is non-linear and as such a correlation between the optical diode and the optical transistor must be established. The optical coupler circuit also has delays related to the operation of the optical diode and the optical transistor, and the operation of the optical diode requires a well defined DC level. As a result, it is generally desirable to avoid the use of an optical coupler circuit.\nA next generation of feedback control does not use optical control circuitry. Instead, the transformer is used to convey real-time feedback signaling from the secondary side to the primary side. In such an application, the transformer includes an auxiliary winding on the primary side that is magnetically coupled to the secondary winding. FIG. 2 illustrates a conventional regulated power converter including a magnetically coupled feedback circuit. The power converter 32 is configured as a traditional flyback type converter. The power converter 32 includes an isolation transformer 34 having a primary winding P1 and a secondary winding S1. The primary winding P1 is electrically coupled to an input voltage Vin and a driving circuit including a transistor 44, a resistor 46, and a controller 42. A capacitor 58 is coupled across the input Vin and coupled with the primary winding P1. Input voltage to the circuit may be unregulated DC voltage derived from an AC supply after rectification and filtering. Similar to the power converter in FIG. 1, the transistor 44 is a fast-switching device controlled by the fast dynamic controller 42 to maintain a desired output voltage Vout. The secondary winding voltage is rectified and filtered using the diode 36 and the capacitor 38, with the output voltage Vout delivered to the load 40.\nThe power converter 32 has a feedback loop that includes a magnetically coupled feedback circuit coupled to the secondary winding S1 of the transformer 34 and the controller 42. The magnetically coupled feedback circuit includes a diode 48, a capacitor 50, resistors 52 and 54 and an auxiliary winding 56. The auxiliary winding 56 is coupled in parallel to the series of resistors 52 and 54.\nThe voltage VA is proportional to the voltage across the auxiliary winding 56. The voltage VA is provided as a feedback voltage VFB to the controller 42. The current through the transistor 44 is also provided as feedback current IFB to the controller 42. The controller 42 includes a real-time waveform analyzer that analyzes input feedback signals, such as the feedback voltage VFB and the feedback current IFB.\nThe auxiliary winding 56 is also magnetically coupled to the secondary winding S1. When the current through the diode 36 is zero, the voltage across the secondary winding S1 is equal to the voltage across the auxiliary winding 56. This relationship provides means for communicating the output voltage Vout as feedback to the primary side of the circuit. The voltage across the auxiliary winding 56 is measured when it is determined that the current through the diode 36 is zero, which provides a measure of the voltage across the secondary winding S1 and therefore the output voltage Vout.\nThe feedback voltage VFB when the diode 36 current is zero is determined and is referred to as the “voltage sense”, and the feedback current IFB when the diode 36 current is zero is determined and is referred to as the “current sense”. The voltage sense and the current sense are input to the real-time waveform analyzer within the controller 42. FIG. 3 illustrates a functional block diagram of a conventional real-time waveform analyzer 60. Error amplifiers 62 and 64 are acceptors of the regulating means, which in this configuration are the sensed voltage and the sensed current. The error amplifier compares the input sensed voltage to a reference voltage and outputs a first difference value. The first difference value is amplified by the gain of the error amplifier 62. The error amplifier 64 compares the amplified first difference value to the sensed current and outputs a second difference value that is either High or Low. A pulse width modulation (PWM) block 66 is configured as a Flip-Flop digital device. The output of the PWM block 66 is set according to the switching frequency of the clock 68 and is Reset by the High or Low value input from the error amplifier 64. The variable signal applied to the Reset pin generates an output signal that is a pulse train modulated by the pulse width. A multiple input OR gate 70 inputs the clock signal, the pulse train signal, a shutdown signal, and a OVP/UVP/OTP signal, where OVP stands for “over voltage protection”, UVP stands for “undervoltage protection” and OTP stands for “over temperature protection”. The waveform analyzer is configured to output a high voltage value when one of the inputs to the OR gate is high or to output a low voltage value when all of the inputs to the OR gate are low. The high voltage value output from the waveform analyzer corresponds to turning on the transistor 44 in FIG. 2. The low voltage value corresponds to turning off the transistor 44. The OR gate also enables a high voltage signal output from the PWM block 66 to propagate to the output by monitoring abnormal conditions such as under voltage, over voltage, over temperature, etc. In this manner, the pulse width of each pulse can be modified which adjusts the output voltage into regulation.\nIn general, control intricacies of the waveform analyzer are aligned with control argument sampling to achieve overall system functional performance. Sampling argument is in the form of current, voltage and impedance. System functional performance is in the form of pulse width modulation (PWM), pulse frequency modulation (PFM) and pulse amplitude modulation (PAM). The waveform analyzer of FIG. 3 is limited to signal processing in DC or real-time switching waveforms. In either case, the feedback signal received by the waveform analyzer requires some status integrity, such as no noise on the DC level, no disturbance on the switching waveform and to some degree represent a combination of analog and digital representations. The voltage across the auxiliary winding typically forms a pulse train with frequency corresponding to the switching frequency of the driving transistor. The voltage across the auxiliary winding when the secondary winding current is zero, which corresponds to the diode 36 current equaling zero, corresponds to the falling edge of the pulse. As such, measuring an accurate voltage value requires that the pulse is well defined with sufficient pulse integrity particularly at the falling edge. Further, the voltage value immediately following the rising edge includes ringing due to the leakage impedance of the transformer. As such, pulse integrity also requires sufficient time for the voltage value to stabilize following the rising edge. Higher switching frequencies minimize the pulse width and therefore provide less time for voltage stabilization. For at least these reasons, providing a pulse with sufficient pulse integrity is often difficult to achieve.\nAnother type of feedback mechanism uses a second transformer as a gate drive transformer controlled by a pulse width modulator on the secondary side of the circuit. FIG. 8 illustrates a conventional regulated power converter including a second transformer used as a gate drive transformer. The power converter 74 is configured as a traditional flyback type converter. The power converter 74 includes an isolation transformer 94 having a primary winding P1 and a secondary winding S1. The primary winding P1 is electrically coupled to an input voltage Vin and a driving circuit including a transistor 80, a resistor 78, and a diode 82. A capacitor 84 is coupled across the input Vin and coupled with the primary winding P1. Input voltage to the circuit may be unregulated DC voltage derived from an AC supply after rectification and filtering. The driving circuit also includes a diode 96, a controller 86, and a gate drive transformer 76 having a primary winding P2 and a secondary winding S2. The diode 96 and the controller 86 are coupled in series to the secondary winding S2 on the secondary side of the circuit. The primary winding P2 of the transformer 76 is coupled to the gate of the transistor 80 via the diode 82. The transistor 80 is a fast-switching device controlled by the fast dynamic controller 86 to maintain a desired output voltage Vout. The voltage across the secondary winding S1 of the transformer 94 is rectified and filtered using the diode 92 and the capacitor 90, with the output voltage Vout delivered to the load 80 and to the controller 86. The secondary side controller 86 regulates the output voltage Vout by transmitting a driving signal to the transistor 80 via the gate drive transformer 76."}
{"text": "This invention relates to the jointing of armoured submarine cables and, more particularly, to the jointing of the armour of a submarine optical fibre cable.\nFor the understanding of this invention, it is convenient to regard an armoured submarine cable as a core surrounded by armour. The core includes transmission elements, which may be optical fibres, strengthening members, insulation and a water seal capable of withstanding the pressure at the depth of operation. Thus, the core, which is complete in itself, provides all the functions needed for transmission.\nWhen it is used in shallow water the cable may be subject to tensile stress; for example, if it is fouled by anchors. Such tensile stress can be in the region of 70 tons. It is important that tensile stress is not allowed to lead to elongation of the core. In the case of cores whose transmission elements are optical fibres this is particularly important because of their high susceptibility to damage by elongation. It is conventional to armour the core so that it is better able to resist tension. The armour usually comprises single or multi start helices of tensile wires wound round the core. Most often the armour comprises one or two such layers. If the cable is pulled, the tension is taken by the armour and relatively little, if any, is transmitted to the core.\nSince submarine cables extend for many hundreds of kilometers there are necessarily many joints. Each joint is conveniently considered as two elements, i.e., the core-joint through which transmission occurs and the armour-joint which transmits tension and protects the core-joint from deleterious effects of tension. The core-joint is made in accordance with the technology of the core and it is usually enclosed in an injection moulded polyethylene envelope. This invention is not concerned with the core-joint.\nKnown armour-joints include the following:\n(i) a cone splice wherein the armour wires are brought through a hollow cone and bent back over its outer surface;\n(ii) a barrel splice wherein the armour wires are passed through a common sleeve from opposing directions and bent back over its outer surface;\n(iii) a ring splice wherein each armour wire is passed individually through a hole in a common ring, being secured therein by a terminal ferrule; and\n(iv) an overlay splice wherein a first set of armour wires is wound helically as an outer layer over the second set of armour wires as an inner layer for an overlap distance of, typically, 8 m, the effect of tension in use being to reduce the internal diameter of the outer layer so clamping it onto the inner layer.\nBoth the cone and the barrel splices have been found to suffer from a lack of strength where the armour wires are bent back. Such joints will carry a significantly reduced tensile load in comparision with unjointed armour wires. The ring splice, in practice, has also been found to suffer from a lack of strength because of the difficulty of installing the individual terminal ferrules at exactly equivalent positions at the ends of each of the armour wires. Any tensile load may not be evenly distributed, as a result, between the armour wires. Lastly, the overlay splice is particularly subject to elongation under load which is inseparable from the reduction in internal diameter. Elongation of the armour-joint can clearly affect the proportion of the tensile load transferred to the core and can lead to elongation of the core which, as mentioned above, can be extremely undesirable."}
{"text": "A fluorescent sensor can produce a detectable change in fluorescence response upon interacting with an analyte. Fluorescent sensors can provide desirable properties such as water solubility, low detection limits, and high selectivity for a desired analyte. The analyte can be a small molecule or ion (such as, for example, H+ (i.e., pH), Ca2+, glucose, or O2)."}
{"text": "In the supply of ink into the nozzle head in an ink jet recorder, it is necessary that pressure of the ink be maintained constant and air bubbles be not mixed in the ink. As a matter of fact, however, it is said that the concentration of air (oxygen) dissolved in the ink in the state that the ink is only exposed to the air is increased to 5 to 6 ppm from 1 ppm for a period of 30 minutes. And, the entrained air accumulates at such as a bent area in the piping to become air bubbles which in the nozzle head cause an undischarge of ink or a fluctuation in ink pressure which in turn brings about a poor or defective printing.\nAccordingly, the ink furnishing unit of this type in the prior art includes a deaerator that is provided for a recovery tank which reserves the ink recovered from the purge section and abounding with air bubble (see, for example, JP P2003-127435 A) or a deaerator that is disposed midway in the ink supply circuit (see, for example, JP P2004-17517 A).\nIn the technique shown in JP P2003-127435 A and designed to permit recycling the ink recovered at the purge section, it is difficult to deaerate the entire recovery tank uniformly so that it has not been possible to remove air bubbles completely from the ink in the recovery tank.\nWhile in this regard the technique shown in JP P2004-17517 A permits the entire volume of ink furnished from the supply tank into the main tank to be evenly deaerated only by passing through its piping, this techniques does not allow such ink to be directly furnished to the nozzle head. Since the ink is once stored in the air-open main tank for stabilization of its supply pressure and then is furnished to the nozzle head, the amount of air entrained in the ink while in the main tank increases with the lapse of time, causing the problem that the ink discharge at the nozzle head tends to be defective.\nWith these taken into account, it is an object of the present invention to provide an ink furnishing unit in an ink jet recorder that can efficiently deaerate the ink furnished to the nozzle head including the ink recovered from the purge section so that the defective printing due to the mixing of air bubbles in the ink may be prevented."}
{"text": "Normally, in continuous casting plants, the molten metal produced, for example from an electric furnace, or from a converter, is subjected to a series of pouring passages between containers, prior to being cast in the ingot mould to form ingots, slabs, bars, strips and others.\nA typical example is given by the application of the technology of continuous strip casting in which the steel coming from the production oven is collected in a ladle, from which it is unloaded into one or more tundishes, each of which feeds one or more melts above two cooled, counter rotating crystallising rolls, which constitute the ingot mould in which the metal solidifies and exits in the form of finished product, for example a strip, a bar or otherwise. Between the tundish and the ingot mould there can also be provided an under tundish, or another discharging device, from which the liquid metal is dosed into the ingot mould.\nIt is known that the quality of the final product, and also the execution of the process itself can be compromised by the contact of air with the molten metal. For example, the oxygen in the air can combine with elements dissolved in the steel forming inclusions, which deteriorate the quality of the steel, as also with dissolved oxygen itself. In case the steel is used in continuous casting, for example in a dual roller plant, the oxygen combined with the iron can produce scale which deposits on the rollers, locally altering, amongst others, the heat exchange, with serious consequences for the final product. Nitrogen can also provoke the formation of precipitates which compromise the quality of the product.\nIn the containers into which the molten metal passes, there is generally maintained a protective atmosphere, generally through the feeding of an inert gas, for example argon.\nThe various passages of the metal from one container to another are critical, with regard to the problem of possible contamination with air, bearing also in mind that the various containers must be emptied in between, both for the normal movement during the metal production operation, as for example in the case of the tundishes, and to allow for the substitution of the parts in refractory material, such as, for example the discharging sleeves or tubes through which the molten metal flows, subject to rapid wear and corrosion.\nDischarging devices in the known art are not easy to operate and require delicate operations to allow the approaching of two containers for joining. Sometimes the impacts produced during the coupling damage the parts of the discharging device. Therefore, the need is felt of providing discharging devices, to interpose between the various containers in which the molten metal passes, avoiding contact between the air and the metal and allowing a rapid and precise coupling between the containers, in addition to a likewise rapid uncoupling."}
{"text": "The present invention is related to a rock climbing apparatus, and more particularly to a lightweight rock climbing apparatus which has sufficient strength and can be easily assembled for personal or domestic rock climbing exercise.\nRock climbing has become a popular activity. However, natural outdoor rock field can be hardly found. Moreover, there is safety problem in real rock climbing. Therefore, many artificial rock fields have been built. However, the artificial rock field necessitates larger space and the equipments of the artificial rock field are expensive. Therefore, the cost for such rock climbing exercise is high.\nU.S. Pat. Nos. 5,549,195, 5,125,877 and 5,919,117 disclose rock climbing apparatuses. In each of those rock climbing apparatuses, a caterpillar structure is arranged on a frame body. The caterpillar structure is connected with a transmission mechanism for restricting the rotational speed of the caterpillar structure. Two sides of the caterpillar structure are respectively disposed in the rails of the frame body via chains or rollers, whereby the caterpillar structure can circularly revolve along the rails. Multiple rocks are arranged on the caterpillar structure for a user to climb for rock climbing exercise.\nIn the caterpillar structure of the above rock climbing apparatus, simply two sides of the caterpillar structure bear the force, while the middle of the caterpillar structure totally lacks support structure. The rocks are positioned on the middle of the caterpillar structure. Therefore, when a user climbs the rocks, the application force will totally fall onto the middle of the caterpillar structure. In order to avoid swing or deformation of the caterpillar structure, the caterpillar structure is designed with quite strong structure for bearing all the application force in climbing.\nFor example, U.S. Pat. No. 5,549,195 discloses a caterpillar structure composed of multiple solid planks. U.S. Pat. No. 5,125,877 discloses a caterpillar structure composed of multiple metal frames pivotally connected with each other. A board body is inlaid in each metal frame. U.S. Pat. No. 5,919,117 discloses a caterpillar structure composed of multiple extruded board bodies pivotally connected with each other. In order to avoid deformation of the caterpillar structures, in all the above rock climbing apparatuses, the caterpillar structures are designed with quite strong and complicated structures. This results in that the rock climbing apparatus has a weight up to several hundred kilograms. Therefore, it is hard to move the rock climbing apparatus. In other words, such rock climbing apparatus is unsuitable for personal or domestic rock climbing exercise.\nFurthermore, the frame body of the rock climbing apparatus is a huge and heavy structure for supporting the heavy caterpillar. Therefore, the rock climbing apparatus can be hardly assembled or transferred.\nBesides, the caterpillar structure has gaps between the board bodies and the frames. These gaps are enlarged when the caterpillar is curved. A user's hairs, clothes, trousers or shoelaces may be chucked in the gaps to result in accident.\nU.S. Pat. No. 6,860,836 of this applicant discloses a relatively lightweight rock climbing apparatus. A circularly rotating mechanism is mounted between two frames of a bed. The circularly rotating mechanism includes several sprockets. Each chain has multiple lugs arranged at intervals for mounting multiple transverse beams in parallel to the shaft. Two ends of each transverse beam are respectively inserted in open sides of the two frames. The transverse beams connected with the chains can be circularly revolved along the inner edges of the open sides of the frames. The chains and the transverse beams are wrapped with serially connected soft pads for shading the chains and the transverse beams. Multiple rocks are locked on outer sides of the soft pads corresponding to the lugs of the chains. Several U-shaped stopper boards are fixed between the frames corresponding to the chains, whereby the chains can lean on the stopper boards.\nThe chains can circularly revolve in a fixed path between the frames. The rocks are arranged on the chains which serve as the main supports for the rocks. Multiple transverse beams are arranged on the chains at intervals for bearing the outward pulling force exerted onto the rocks. The chains lean on the U-shaped seats for bearing the inward pushing force exerted onto the rocks. The above structure is simplified and the weight thereof is minified.\nHowever, in practice, the above rock climbing apparatus still has some shortcomings. The rocks are fixed on the chains. When the rocks are forced inward, the lateral sides of two links of the chain abut against the U-shaped seat. When the chain revolves, the links and the U-shaped seat will abrade each other to produce noise. After a period of use, the links and the U-shaped seat will be worn out. This will affect the smoothness of the revolution and even the strength of the chain. In addition, when the rocks are forced outward, two ends of the transverse beam will abrade the open sides of the frames. Similarly, the transverse beam and the inner sides of the frames will be worn out and noise is produced.\nMoreover, the above rock climbing apparatus is lightweight and suitable for personal or domestic use. However, the frames still have a considerable length exceeding the height of an adult for a user climb. The frames are divided into two segments. However, the segments can be hardly conveniently DIY assembled. Furthermore, there is no support structure between the frames. The entire structure is simply supported by the bed. Therefore, the strength of the rock climbing apparatus is still insufficient.\nBesides, the transverse beams are arranged at intervals. Therefore, a gap is defined between two transverse beams. Although the transverse beams are wrapped with the soft pads, a user still often touches the gaps when climbing. At this time, the soft pads will be depressed to contact with the chains or the skeleton inside the frame. Under such circumstance, the inner faces of the soft pads will be abraded and worn out by the chains or the skeleton of the frame.\nIn addition, the soft pads must have a certain tightness to keep the surface smooth. However, the transverse beam is wrapped with the soft pad and fixed on the chain via the rock. When the transverse beam is revolved to the sprocket, the transverse beam will be tangential to the circumference of the sprocket. Accordingly, two ends of the transverse beam will lift the soft pad as shown in FIG. 12. At this time, the soft pad will be tensioned. In some cases, the soft pad will be broken in the position where the rock is fixed. In some serious cases, the chain will be unable to revolve. In the case that the soft pad is made of a resilient material, the inner face of the soft pad will be abraded by the chain or the skeleton of the frame. As a result, after a period of use, the soft pad will elastically fail due to the continuous tension."}
{"text": "In semiconductor manufacturing facilities, a reactor for generating moisture with a structure as shown in FIG. 15 has been widely used. Namely, in FIG. 15, the following character references are used:    A Reactor    B Temperature regulator    H2 Hydrogen gas    O2 Oxygen gas    G Mixed gas    W Moisture gas    L Gap    V Interior space    α Inclination angle of the peripheral edge of a reflector    1 Reactor structural component on the inlet side    1a Gas feed port    2 Reactor structural component on the outlet side    2a Moisture gas take-out port    3a, 3b Reflectors    4 Reflector fixing bolt    5 Spacer    6 Platinum coating catalyst layer comprising a barrier coat 7 and a platinum coat 8 mounted on its outer surface    9 Barrier coat mounted on the inner surface of the reactor structural component on the inlet side    10 Barrier coat mounted on the outer surface of the reflector    11 Welded portion    12 Fixing hole for the sheath-type thermometer    13 Heater    15 Heater press    15 Cooler    15a Heat dissipation fin    15b Substrate\nA mixed gas G produced by mixing hydrogen gas H2 and oxygen gas O2 is fed from the gas feed port 1a to the interior space V, and when the gases contact with a platinum coating catalyst layer 6 mounted on the inner surface of the reactor structural component on the outlet side 2, the reactivity of the hydrogen and oxygen is activated by catalytic action. The activated hydrogen and oxygen react at the appropriate speed without causing an explosive combustion reaction within the atmosphere, below the combustion temperature of hydrogen, and the high purity moisture gas V produced flows out from the moisture take-out port 2a. \nIt may be necessary to raise the temperature of the interior space V of the reactor A higher than at least 200° C. in order that the afore-mentioned reactivity of the hydrogen and oxygen is activated and to maintain the stable reaction of both hydrogen and oxygen. To achieve this end, a plane-shaped heater 13 is provided on the outer side of the reactor structural component on the outlet side 2, and the reactor A is heated by the plane-shaped heater when the reactor A starts up. FIG. 16 shows the relationship between the temperature of the reactor A and the reaction rate of the hydrogen.oxygen. When the temperature of the reactor exceeds approximately 200° C., the reaction rate of the hydrogen.oxygen becomes the value of approximately 98% or more regardless of the mixture ratio of both hydrogen and oxygen.\nWhile the reaction of the afore-mentioned hydrogen and oxygen proceeds, the reactor A is heated by the reaction heat, and then the reactor's temperature goes up. In order to restrain the afore-mentioned explosive combustion reaction of the hydrogen and oxygen, it is necessary that the temperature of the interior space V is held at a temperature lower (for example, 400˜450° C.) than the lowest limit ignition temperature (approximately 560° C., and the limit ignition temperature goes up higher than 560° C. depending on the mixing ratio of hydrogen and oxygen) of the hydrogen gas (or hydrogen containing gas).\nAccordingly, with this kind of conventional reactor for generating moisture, such measures as limiting the flow rate (namely, the volume of moisture gas W generated), improving the cooling performance of a cooler 15, or increasing the heat capacity of a reactor A, are implemented. When such measures as these are employed, it becomes possible to produce high purity moisture gas at the low cost, and stably and continuously with this kind of reactor for generating moisture. Thus, such reactors are widely and practically used.\nHowever, in recent years there is seen a trend to downsize the apparatus for generating moisture while providing a bigger volume of moisture generation, which has been demanded in the field of semiconductor manufacturing facilities. For example, when the bore of a wafer becomes larger, the volume of moisture required for each treatment process increases. Accordingly, there has been a strong demand for more volume of moisture generation than that presently obtained by a conventional moisture generating apparatus having a dimensional volume similar to that of the conventional ones that are produced.\nHowever, when downsizing of semiconductor manufacturing facilities is required, the moisture generating apparatus has to meet a severe restriction placed upon its volume/dimension. Because of this reason, it is nearly impossible for the cooling performance of the reactor A to be improved by means of upgrading the cooling fan though improvements such as making the heat dissipation fin 15a of the cooler 15 higher and increasing the number of heat dissipation fins slightly. Similarly, it is very difficult to upgrade a reactor A in view of the afore-mentioned restrictions on its volume. For example, the structurally allowable dimensions of a reactor A for generating moisture having a large flow rate of more than 10 SLM (“Standard Liter per Minute”) are similar to the values of dimensions (outer diameter 228 mm, thickness 37 mm and height of a heat dissipation fin 25 mm) of the conventional reactor for generating moisture which has a maximum volume of moisture generation of 5 SLM.\nTherefore, for a conventional reactor equipped with a cooler 15 provided with the structure shown in FIG. 15, there is no way at all to meet the requirement of increasing the volume of moisture generated in view of the restrictions in its volume/dimension including that volume/dimension of the cooler 15. Thus, a conventional reactor equipped with a cooler 15 fails to increase the volume of moisture that may be generated safely and easily.\nPatent Document 1: TOKU-KAI No. 2001-48501\nPatent Document: 2: INTERNATIONA DISCLOSURE WO-01/94254A"}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus of inserting insert or enclosure material such as sheets of document, booklet and the like into an envelope which is made of paper, plastic film or any other appropriate materials. More specifically, the invention relates to envelope inserting method and apparatus which makes possible improvement in the envelope inserting efficiency by permitting the insert material to flow in a continuous, streamlined manner.\nEnvelopes are available in various shapes and sizes, including oblong-shaped envelopes having a flap portion on the longer side or the shorter side thereof, envelops with predetermined sizes prescribed by the post office and various sizes determined according to the content document sizes such as international A5, A4 and A3 sizes, and various DL envelopes with an outlook window.\nThe present invention is applicable not only to such envelopes of various shapes and sizes, but also to various forms of contents to be inserted or enclosed in the envelopes, including a plurality of cut sheets laid in a bundle, insert material folded in half, third or fourth, brochure or books saddle stitched or perfect bound and collated insert material. It is noted that in this specification these enclosure contents to be inserted into envelopes are generally referred to as xe2x80x9cinsert material.xe2x80x9d As is apparent in the art, after the insert material has been inserted, the envelope is closed and sealed by its flap. The present invention is directed to the method and apparatus of envelope inserting which is performed preceding such flap closing and sealing process.\n2. Description of the Related Art\nFor better understanding to the background underlying the present invention, reference is made to FIG. 2 which shows a typical known method of envelope inserting. As seen in the drawing, insert material 1 is being inserted as indicated by arrow 3 into an envelope 2 which is then located at envelope inserting position through an opened throat of the envelope. A preceding envelope 2a having the insert material 1 already inserted therein is transferred in the direction of arrow 3a and a succeeding empty envelope 2b is transferred toward the envelope inserting position as indicated by arrow 3b. \nAs would be appreciated from the drawing, the conventional envelope inserting method is characterized in that all operations involved with the envelope inserting, including insertion of the insert material 1 and the transferring of envelopes 2, 2a, 2b, are performed in a horizontal plane. Another feature of the prior art method resides in the change of moving direction of the insert material 1. That is, the insert material 1 is moved in the arrow direction 3 to be inserted into the envelope 2, while the insert material 1 thus inserted into the envelope is then moved in the direction 3a that is perpendicular to the inserting direction 3.\nThe distance for which the insert material 1 must be moved for insertion into the envelope 2 generally depends on the size of the envelope as measured in the arrow direction 3. Accordingly, longer time is required for the insert material 1 to be inserted completely with an increase of the above size of the envelope. An attempt may be made to design an apparatus which can speed up the envelope inserting operation, although this will only invite complication and hence higher cost of the apparatus.\nOn the other hand, the change of moving direction of the insert material 1 means that insert material 1 and the envelop receiving such material must make a temporary stop before the insert material 1 changes the direction of its movement. If a number of successive sets of insert material 1 is fed continuously toward the envelope inserting position at a predetermined speed, any two successive or adjacent sets of insert material 1 must be fed at a sufficiently spaced interval. This will inevitably limit the envelope inserting operation speed, thus making it difficult to achieve today\"\"s need for high-speed envelope stuffing operation.\nTherefore, it is an object of the present invention to provide a method of envelope inserting which is operable at a rate exceeding what has been achievable by conventional methods by allowing the envelopes and the insert material to move in a rational manner.\nIt is another object of the invention to provide an apparatus for practicing the this envelope inserting method.\nTo achieve the objects of the invention, the invention provides an envelope inserting method wherein the insert material is conveyed along a straight path successively to an envelope inserting position where the insert material is inserted into the envelope and the stuffed envelope is then conveyed out of the inserting position in the same direction as the above straight path. For this purpose, the method of the invention includes providing immediately above the envelope inserting position an envelope storage station where a plurality of envelopes are stored in a stack one above another with the front panel of each of the envelopes facing upward so that the envelopes in the storage station are fed successively to the envelope inserting position for insertion thereinto of insert material.\nIn the preferred embodiment of the invention, feeding of the envelopes is performed by means of a vertically disposed rotatable screw-like spiral member located at the downstream side of the envelope storage station as seen in the direction in which the insert material is conveyed. The spiral member has formed in the periphery thereof a spiral groove for receiving therein the bottom edges of the respective envelopes, thereby holding such envelops at the bottom edges thereof such that the bottom edges are separated one from another. The spiral member is rotatable in such a direction that causes the bottom edges of the envelopes to move downwards with the rotation of the spiral member and then to be released successively from the holding by the spiral member.\nAccording to this method, each empty envelope has to move only for an extremely short distance to the envelope inserting position and the insert material does not have to change its direction of movement after it has been inserted into an envelope, with the result that the flow of the insert material can be streamlined and very smooth and, therefore, a remarkable speedup and improvement in the envelope inserting efficiency can be achieved.\nThe spiral member as the envelope feeding means can also serve to hold the envelope at its bottom edge in the envelope storage station, so that the lowermost envelope in the storage station which is subjected to frictional contact with the stuffed envelope then moving away from the envelope inserting position can be prevented from being dragged and dislocated by the movement of the stuffed envelope.\nThe present invention also provides an envelope inserting apparatus which comprises conveyer means for conveying the insert material along a straight path toward an envelope inserting position of the apparatus where the insert material is inserted by the conveyer means into the envelope through an opened throat and then conveying the stuffed envelope out of the envelope inserting position along the same straight path, an envelope storage station provided immediately above the envelope inserting position for storing a plurality of empty envelopes in a stack one above another with the front panel of each of the envelopes facing upwards, and throat opening means disposed adjacent to the upstream side of the envelope storage station, as viewed in the direction in which the insert material is conveyed, and operable to be brought into contact with the throat portion of envelope for opening the throat of that envelope. The envelope storage station is defined on the upstream side thereof by a flap guide member for guiding the flap portions of the envelopes and on the opposite downstream side by envelope feed means for feeding the envelopes successively to the envelope inserting position.\nThe conveyer means include a pushing member which is adapted to be brought into engagement with the trailing end of the insert material to push the material forward into the envelope and then into engagement with the throat portion of the rear panel of the envelope thereby to push the envelope out of the envelope inserting position along the above straight path. The envelope feed means is provided by a vertically disposed rotatable spiral member of the same structure and operable in the same manner as described in the above with reference to the envelope inserting method. The envelope inserting apparatus thus constructed can offer the above-described features and advantages.\nThus, the envelope inserting method and apparatus according to the invention can make possible remarkable improvement in working efficiency, thereby contributing greatly to fulfillment of today\"\"s demand for further speedup in operation and the economy of the related industry.\nThe above and other objects, features and advantages of the invention will become apparent to those skilled in the art from the following description of preferred embodiments of the envelope inserting method and apparatus according to the present invention, which description will be made with reference to the accompanying drawings, wherein:"}
{"text": "The present invention relates to methods and apparatus for remotely verifying the identity of a person, such as over an intranet, internet or for similar circumstances in which the identity of an individual must be verified.\nThere are many methods and techniques used to verify the identity of a person during a face-to-face transaction. One method is to compare the person\"\"s facial characteristics with an authenticated picture on a validated document. For instance, many businesses require the identity of a patron presenting a check to by validated by a driver\"\"s license which includes a photograph of the individual. Similarly, passports and identification cards issued by state or federal agencies are used for verification of the identity of the individual.\nAnother popular method of verifying the identity of an individual during a face-to-face transaction is through written signature comparison. Many of the previously mentioned documents such as driver\"\"s licenses, passports and identification cards also include a written signature of the holder. By comparing the signature on these identifying documents with a written signature signed in the authenticator\"\"s presence, the identity of the person may be verified as represented by the document. Other methods of verifying a person\"\"s identity are through other biometrics such as fingerprints, distinctive patterns in the retina of the eye, identifying characteristics of hands, and the manner in which an individual grasps a writing instrument. Each of these methods requires specialized hardware or that the authenticator possesses specialized knowledge. While each of these methods can be used to authenticate the identity of a person for a face-to-face transaction, other methods are required to verify a person\"\"s identity for remote transactions. Remote transactions for this purpose are defined to be those transactions in which the individuals conducting the transactions are not face-to-face. Examples of remote transactions include transactions over the telephone, computer transactions over an internet or intranet, and similar transactions.\nAuthentication of an individual in remote transactions typically relies on information purportedly unknown to anyone besides the individual whose identity is being authenticated, e.g., a secret. Information used to authenticate the identity of an individual may include the maiden name of the mother of the individual, social security number, address, zip code, account information such as an AMERICAN EXPRESS(copyright), MASTERCARD(copyright), or VISA(copyright), home telephone number, or unique personal identification number (PIN). With the advent of the Internet and the information available over the Internet, many of these previously acceptable methods of identification are becoming obsolete. Once public information such as the maiden name of the mother of the individual, social security numbers, addresses, zip codes, telephone numbers, and similar information may be accessed over the Internet, this information can no longer be used to verify the identity of an individual.\nTwo methods of determining an individual\"\"s identity without relying on the use of publicly available information for remote transactions include electronic signatures and PINs. For instance, if an individual wanted to access their bank account information via the Internet, the bank may require an electronic signature be sent to the bank to be used to authenticate the identity of the person performing the transaction. Heretofore signatures for remote transactions required specialized hardware on the computer system at which the user performs the transaction; i.e., attempts by users to provide repeatable signature using a mouse or similar pointing device have been unsuccessful. Instead, an electronic pen and a tablet or similar specialized hardware have seen required at the computer terminal being used for the transaction by the individual. Even this method of verification of the identity of a user is not foolproof because electronic signatures may be stored and sent when requested by an individual trying to impersonate another individual. The storage of the electronic signature may be performed on the computer used by the individual for the transaction or by other machines or computers which intercept the electronic signature as it is communicated from the user to, in this example, the bank or other financial institution.\nValidating the identity of an individual by using PINs suffer similar drawbacks. When individuals choose personal identification numbers they typically use small number of predictable numbers correspondence to, for example, their birth date, portions of their Social Security number, a loved one\"\"s birth date, an anniversary date, a date of a significant event to that individual, or some other easily remembered number. While the selection of these numbers or groups of numbers allow the individual to recall the PIN easily, the association of the numbers with the individual reduces the associated security. Many systems requiring authentication attempt to minimize access by impersonators by limiting the number of attempts accepted from an individual. Many automatic teller machines (ATM) are programmed to confiscate bank cards after three attempts to access the information with incorrect PINs have been attempted. Additional problems with PINs are present when PINs are retained in computer systems. Electronically transmitted PINs may also be obtained by unauthorized individuals and used for later access attempts.\nFor many transaction over an internet or intranet authenticating an individual\"\"s identity is also important. Examples of instances in which the identity of an individual should be authenticated include online banking, purchases via the Internet or an intranet, access to non-public databases or information, access to medical records, remote access to computer systems, e-commerce, e-banking, business-to-business (B2B) transactions, business-to-consumer (B2C) transactions, e-learning, e-training and similar circumstances. In an effort to minimize the amount of hardware require at the computer system used by the individual to access the secured information, many of these sites required the individual being authenticated to enter a PIN as their access code.\nOne method to alleviate concerns with PINs includes the generation of a random number for the PIN used to authenticate the individual. Random numbers assigned to individuals create other problems in that individuals typically have a difficult time remembering these random PINs. To allow access to these protected systems these individuals may write these PINs down, and store them in their daily planners, on their calendars, in their wallets or purses, or similar locations. By writing these randomly generated PINs down the security of the overall system may be compromised and another individual may gain access to these protected systems by obtaining the recorded PINs.\nWith increasing computer power, methods incorporating electronic signature or digital signature recognition and authentication systems have been advanced. Such systems typically include an input device such as a digitizing pad or tablet to capture and digitally sample the written signature image and/or a biometric feature of the written signature in various ways to compare the new signature to a previously-stored xe2x80x9cauthenticxe2x80x9d exemplar signature. Currently, written signature authentication solutions fail to provide an effective and particularly reliable signature authentication/verification system which may be readily commercially implemented. Furthermore, with the increasing use of the Internet for a myriad of applications and transactions, accurately and reliably verifying a signature on-line is particularly desirable.\nAn object of the present invention is to provide a system and method of validating the identity of a user over an intranet or the Internet which does not require specialized hardware at the computer system used by the individual whose identity is to be validated.\nA further object of the present invention is to provide a method and apparatus that preferably employs a simple method of creating a user unique, i.e., signature using a pointing device, such as a mouse and a Virtual Pad at an End-User/Remote User PC (Personal Computer).\nThe invention is directed to a system and method of authenticating a user by sensing information about the user (biometrics) and eliciting information from the user (a secret). While applicable to the e-commerce environment, the invention is equally usable in a wide range of applications wherein a manually manipulated pointer device such as a computer mouse is available. Using the pointer or pointing device (hereinafter, xe2x80x9cPDxe2x80x9d), the user draws lines and drags (repositions) and/or clicks on icons positioned on a background image to create a user PD signature. The combination of lines, icon repositioning and specifically placed and timed series of clicks are used to identify the user by comparison with a stored PD signature. In addition to the xe2x80x9csecretxe2x80x9d information representing the series of actions using the PD, biometric data is also extracted. The biometric data includes, for example, PD manipulation kinetics (e.g., positioning speed, acceleration, shake, etc.), user click duration, etc. Other biometrics may be added including, for example, a thumbprint reader integral to the PD. Preferably, user profile information, including user identification information, is stored in an encrypted format to protect such information from unauthorized access and disclosure.\nThe invention may use a variety of methods to learn, store and update PD signature profile information to authenticate a user. A preferred method periodically stores sampled PD positional information as vectors including a pair of position coordinate values (e.g., xe2x80x9cxxe2x80x9d and xe2x80x9cyxe2x80x9d values), a sample time (xe2x80x9ctxe2x80x9d) and a PD mode. The PD mode may be in the form of a numeric code representing such conditions as, for example, a normal mode (e.g., no button depressed), a click mode (e.g., a left mouse or equivalent button depressed), and a xe2x80x9cdrag and drop modexe2x80x9d (identifying a particular draggable icon has been captured in response to a mouse click). This stored sampled PD positional information may also be updated with information contained in later authenticated PD signatures.\nAnalysis of the vectors may include identification of nodes or xe2x80x9cagentsxe2x80x9d at predetermined distances along a signature trail of a sampled signature to be learned. Once initially identified and located, the agents are repositioned using an energy or cost minimization technique such as Adaptive Resonance Theory-Fuzzy Cluster Means (ART-FCM.) The repositioned agents then form the starting point for future signature analysis. Thus, in a verification mode, the geometry of the agents of a signature under examination after agent energy or cost minimization is compared with the geometry of the agents for the learned or authenticated exemplar PD signature, the comparison being indicative of a degree of similarity between the two PD signatures. The degree of similarity may be represented by some standard metric, such as a sum of squares of the deviations between geometry parameters. Of course, other matching techniques may be used including, for example, wavelet analysis of the PD signature.\nUse of the invention in the context of e-commerce is particularly advantageous to authenticate a purchaser in an open network environment. For example, the identity of a purchaser may be verified by redirecting a web client from a vendor\"\"s web site to a PD signature server. The server identifies the user (this information being automatically sent by the vendor web site or provided by the client machine) and downloads to the client a predefined virtual pad consisting of a background image and prepositioned draggable (and/or fixed) icons. The prospective purchaser uses a pointing device, such as a standard computer mouse, to draw one or more lines, and click on and/or drag icons to predetermined positions on the background. The series of actions are captured as the vectors including mouse x and y coordinates (with reference to the background image), time and PD (or mouse) mode. The server uses the vectors (together with other global information and security checks) to verify the purchaser\"\"s identity and send a corresponding message to the e-commerce web site so as to authorize (or deny) some action such as a purchase."}
{"text": "1. Field of Technology\nThe present disclosure relates to devices and methods for repairing tissue.\n2. Related Art\nArthroscopic procedures often require soft tissue to be reattached to bone. To achieve this, anchors are placed in the bone and sutures attached to the anchor are passed through the tissue to securely retain the tissue in place. The anchors may be of a one-piece design or may be of a multiple piece design. These anchors may be made from a material that allows the anchor to be absorbed by the body over time. However, suture anchors made from such absorbable material may not allow for the tissue to be held securely against the bone during the degradation of the anchor, thereby not allowing for reattachment of the tissue to the bone. In addition, these anchors don't allow for enough surface contact between the anchor and the device used to insert the anchor into bone, thereby reducing the ability of the user from applying a high amount of torque to the anchor. As a result, the user is forced to create a hole in the bone prior to insertion of the anchor, thereby adding procedural steps and time to the surgery."}
{"text": "In general, an ink composition containing titanium oxide has the problem that it has a high specific gravity and therefore settles down and separates with the passage of time. Accordingly, a means in which the ink composition is sufficiently stirred again before use and then used or a means in which the ink composition is provided with a shear thinning property to inhibit titanium oxide from settling down has so far been taken.\nIt has so far been known as the above means for providing an ink composition with a shear thinning property to add a thickener and a gelatinizing agent thereto. Materials having various chemical structures have so far been investigated for the above additives. However, it is the existing situation that materials investigated for ink compositions (so-called oil-based ink compositions) comprising organic solvents as principal components are few as compared with materials added to ink compositions (so-called water-based ink compositions) comprising water as a principal component.\nKnown as a thickener and a gelatinizing agent which are added to oil-based ink compositions are, for example, montmorillonite base clay minerals, fluorinated phlogopite, dextrin fatty acid esters, fatty acid amides and aluminum 2-ethylhexanoate.\nAddition of the above materials to the intended ink compositions has made it possible to inhibit to the utmost a hard cake from being formed by settling of titanium oxide up to about 2 months. However, there has been the problem that when stored over a long period of 2 months or longer, a settling layer is gradually formed at the bottom of the ink. Further, a gel structure has strongly been formed in the whole part or a part of the ink compositions with the passage of standing time, and reduction in fluidity and a leveling property of the ink compositions following reduction in a shear thinning property has been brought about. On the other hand, considering that the ink composition is provided to the market in the form of a product, it is a matter of course that the quality is required to be maintained over a longer period than 2 months. Accordingly, a coating fluid in which a quality is maintained over a longer period is desired to be provided.\nA material which is an additive investigated in the present invention and which comprises an N-acylamino acid derivative whose fundamental frame is at least one selected from N-acylamino acid amides and N-acylamino acid esters is also known as an oil gelatinizing agent (refer to, for example, patent document 1). Known as an ink composition containing an additive having such characteristic are, for example, one in which 0.5% by weight of N-lauroyl-L-glutamic acid-α,γ-di-n-butylamide is added (refer to, for example, patent documents 2 and 3), one in which 5% by weight of N-lauroyl-L-glutamic acid-α,γ-di-n-butylamide is added (refer to, for example, patent document 4) and one in which 1 to 10% by weight thereof is added (refer to, for example, patent document 5).\nAs described in patent document 1 described above, however, these additives are scarcely soluble in organic solvents, and therefore a specific means in which inorganic metal salts are allowed to coexist as a solubilizing agent is required. Further, it is known that these additives have a strong oil coagulant function and form a three-dimensional network in a state of enclosing oil in an individual net, so that gel which is short of fluidity is formed.\nFurther, as described in patent documents 2 and 3 described above, marked reduction in the fluidity is not observed in an ink composition containing 0.5% by weight of the gelatinizing agent because of a small content of the gelatinizing agent, and the writing characteristic is good. However, there has been the problem that titanium oxide can not be inhibited from settling down over a long period of 2 months or more to cause settling and separation as time passes. In an ink composition in which 5% by weight or more of the gelatinizing agent is added, a net work structure of gel is sufficiently formed because of a large amount of the gelatinizing agent, and titanium oxide can be inhibited from settling down over a long period of 2 months or more. However, fluidity of the ink composition is a little short due to the strong net work structure of the gel, and the gel becomes further firm as time passes, so that observed is the problem that reduction in fluidity and a leveling property of the ink composition is brought about to a large extent.\nOn the other hand, also in the ink composition in which the gelatinizing agent is suitably added, as described in patent document 5 described above, in a range of 1 to 10% by weight, particularly in a range of 1% by weight or more and less than 5% by weight, brought about to not small extent is either of the problems that settling of titanium oxide is caused when time of 2 months or longer passes or that the gel becomes firm to cause reduction in fluidity and a leveling property of the ink composition, and it is the existing situation that both problems (first problem) do not result in being solved at the same time.\nFurther, it has so far been known that an organic solvent having a relatively high volatility is used for a coating fluid in a fluid applicator such as a writing instrument, a correction device, an adhesive and a toilet tool for the purpose of quickly drying the coating fluid when it is coated. However, organic solvents used for paint markers and correction devices have a high volatilizing rate, and therefore caps have to be put thereon in non-use to prevent the coating fluids from drying. Accordingly, involved therein are the problems (second problem) that a labor for removing the caps in every use is required and that if the caps are neglected to be put on after use, the coating fluids are dried and solidified for a short time, so that the applicators can not be used.\nOn the other hand, known are a lot of correction devices of a ballpoint pen type equipped with a means in which a ball made of metal is pressed against a point aperture part, for example, by a pressing means such as a spring to seal it (refer to, for example, patent documents 6 and 7 filed by the present applicant).\nOn the other hand, a lot of ink compositions for marking pens which are prevented from drying up by adding a cap-off performance-improving agent such as lecithin to a coating fluid have so far been known, but if a cap-off performance-improving agent such as lecithin is used for a coating fluid containing titanium oxide which is liable to settle down, coagulating and settling down of titanium oxides themselves caused by breakage of a dispersion system and subsequent separation of the liquid are brought about in a certain case, and therefore it is the existing situation that cap-off performance-improving agents used for ink compositions for writing instruments can not simply be diverted. Further, the present applicants have filed a correction liquid-blended composition characterized by adding at least paraffin waxes as a cap-off performance-improving agent (refer to, for exampled, a patent document 8).\nHowever, in the art described in the patent documents 6 and 7 described above, it is impossible to completely seal a tip part because of contact of metal with metal in which a metal ball is pressed against an aperture part, and the problem (third problem) that a cap is always put on in non-use is not solved.\nFurther, the art described in the patent document 8 described above is a correction liquid composition having an excellent cap-off performance which has not so far been available. However, if an addition amount of paraffin wax is small, a period in which the cap-off performance can be maintained is about 2 weeks, and on the other hand, a period in which the cap-off performance can be maintained is extended to one month or longer by increasing an addition amount of paraffin wax. In the product, however, a quality has to be maintained for a long period of several months or longer, and therefore it is the existing situation that the product can not help being equipped with a cap even when the art described above is used.\nFurther, in a gel ink ballpoint pen for correction which is a fluid applicator, a fluid reservoir filled with a correction liquid (fluid) has so far been desired to be so-called refillable because of an environmental problem. In such case, a ball point has to be protected in transporting and storing a refill alone, and a solvent has to be inhibited from being volatilized from the point.\nIn order to prevent such volatilization, countermeasures such as making parts constituting a fluid reservoir main body of metal or resins and providing a tip point with a seal by a resin coating film have been taken.\nMany arts regarding the above fluid reservoir in which a tip point is sealed by a resin coating film have so far been proposed, and known are, for example, a resin coating film formed from a synthetic resin emulsion (refer to, for exampled, patent document 9), a resin coating film formed from a molten thermoplastic resin (refer to, for exampled, patent document 10), a resin coating film formed from a polymer latex (refer to, for exampled, patent document 11), a resin coating film formed from a hot melt resin (refer to, for exampled, patent document 12) and a colored resin coating film (refer to, for exampled, patent document 13).\nThe seals formed by the resin coating films described in patent documents 9 to 11 and 13 described above prevent solvents having no high volatility and water contained in water-based inks and water-based gel inks having a shear thinning property from vaporizing from tip point parts (writing point parts). In the cases of applicators filled with fluids containing solvents having high volatility such as n-hexane, cyclohexane and methylcyclohexane, the solvents are still insufficiently inhibited from vaporizing even if resin coating films of synthetic resin emulsions and polymer latexes are formed, and therefore involved therein is the problem (fourth problem) that the fluids are liable to be solidified at the sealed parts due to dry-up so that the fluids can not discharge from the tips when long time passes from production through use.\nAlso, the seal prepared by the resin coating film described in patent document 12 described above is formed by coating a tip point of an applicator filled with a coating liquid having a high viscosity such as a correction liquid, a toilet liquid, an adhesive and a paint with a hot melt resin such as an ethylene-vinyl acetate base resin, a polyester base resin, a polyamide base resin, a polyolefin base resin and a rubber base resin, which are comprehensive resin names, and it shows an art close to that of the present invention. However, careful investigation of the above patent document 12 shows that the examples are not described and that the specific kinds of the solvents and the specific names of the hot melt resins are not described.\nIn particular, in the cases of applicators filled with fluids containing solvents having high volatility such as n-hexane, cyclohexane and methylcyclohexane, the solvents are still insufficiently inhibited from vaporizing even by resin coating films of ethylene-vinyl acetate base resins and polyamide base resins, and therefore involved therein is, as is the case with the fluid applicator for a water-based ink described above, the problem (fifth problem) that the fluids are liable to be solidified at the sealed parts due to dry-up so that the fluids can not discharge from the tips when long time passes from production through use.\nPatent document 1: Japanese Patent Application Laid-Open No. 68102/1977 (examples and others)\nPatent document 2: Japanese Patent Application Laid-Open No. 343875/2000 (examples and others)\nPatent document 3: Japanese Patent Application Laid-Open No. 343879/2000 (examples and others)\nPatent document 4: Japanese Patent Application Laid-Open No. 53034/1997 (examples and others)\nPatent document 5: Japanese Patent Application Laid-Open No. 256178/2002 (examples and others)\nPatent document 6: Japanese Patent Application Laid-Open No. 58292/1996 (claims, examples and others)\nPatent document 7: Japanese Patent Application Laid-Open No. 118896/1996 (claims, examples and others)\nPatent document 8: Japanese Patent Application Laid-Open No. 213162/2003 (claims, examples and others)\nPatent document 9: Japanese Patent Application Laid-Open No. 96597/1983 (claims and others)\nPatent document 10: Japanese Patent Application Laid-Open No. 35294/1986 (claims, examples and others)\nPatent document 11: Japanese Patent Application Laid-Open No. 98396/1982 (claims, examples and others)\nPatent document 12: Japanese Patent Application Laid-Open No. 222950/1995 (claims, examples and others)\nPatent document 13: Utility Model Registration No. 3076019 (claims, examples and others)"}
{"text": "1. Field of the Invention\nThe present invention relates to a direction and distance estimation system and method based upon characteristics of the Radio Frequency (RF) and the power path loss curve to locate victims.\n2. Description of Related Art\nAircraft emergency locator systems such as ELTs (Emergency Locator Transmitters) and terrestrial location systems such as PLBs (Personal Locator Beacons) are well known in the existing art. Individuals lost at sea and fortunate enough to have emergency location identification devices such as an Emergency Position Indicating Radio Beacon (EPIRB), can send out RF signals with the hopes that their distress signals will be picked up by maritime stations or satellites, which can in turn, relay the information to the proper authorities.\nHandheld or mobile location devices and systems can effectively locate the directional parameters associated with distress beacons; however, in order to obtain accurate distance readings, one must employ the use of satellites and global positioning systems (GPS). If a distress beacon emits a GPS signal to orbiting satellites, then a searcher may easily locate the distress beacon based upon the GPS reading. Many distress beacons currently being used are not equipped to transmit GPS readings and merely emit a RF signal for reception by a handheld, mobile, stationary or orbiting satellite receiver. As stated above, the RF receivers used to track and locate a RF signal from a distress beacon effectively locate the direction of the beacon, however reliable distance readings under this application are not possible. The lack of distance readings to the beacon can unnecessarily prolong or possibly imperil a search and rescue operation. If an individual is attempting to locate a person lost at sea, then initial movement of a search vessel is carefully monitored in order to prevent possibly running over the distressed victim. If the search vessel is a great distance from the distressed victim, then valuable time may be lost due to the necessity to move slowly during the searching process. If accurate distance readings were available, then the searching process could be shortened and therefore increase the likelihood of a successful rescue.\nDirectional and distance information may be equally important to search and rescue operations on land. Searchers may need to locate a distress beacon in dense woods or an urban metropolitan area. The beacon may be used by hikers, off road adventurers, skiers, campers and children. If the searcher is limited to directional data only, then even the effectiveness of a land search could be significantly diminished.\nAvalanches are a great hazard to skiers, hikers and climbers and other people in mountain regions during winter. If a person is buried by snow, there is a very short period of time available to recover the person before he or she suffocates or dies of exposure due to the snow pack. The time factor in locating someone buried in an avalanche is a critical factor in the survival of the person buried. Electronic locating devices are known in the prior art to find victims of an avalanche that are buried under snow. Avalanche rescue devices manufactured by several companies transmit an RF signal on a frequency of 457 kHz. U.S. Pat. No. 6,031,482 issued to Lemaitre, et al., Feb. 29, 2000, discusses a method and system for sensing and locating a person, e.g. under an avalanche. Various types of electronic systems are discussed in this patent. U.S. Pat. No. 6,167,249 issued to Hereford, et al. discloses an avalanche victim locating transceiving apparatus that uses a plurality of different antennas that are arranged to provide directional and distance estimates to the victim along the lines of flux. None of the references exploit the RF characteristics of the field or power path loss curve to estimate a point to point direction and distance to an avalanche victim.\nIt would therefore be advantageous to provide a method and system for the location of a distress beacon that would provide the searcher with accurate distance and direction estimations. By providing a searcher with distance and directional estimations, the search time may be reduced significantly without endangering the distress victim. Although, the search and rescue operation has been discussed in association with sea and water rescue, distance estimations are also advantageous in regard to terrestrial search and rescue missions. It would also be advantageous if a handheld device facilitated the operation of the present system and method."}
{"text": "Various embodiments of electrical plug-in connectors are known from the prior art. Electrical plug-in connectors have for example a first and a second contact casing with contacts, the two contact casings being able to be pulled by means of a pivotable stirrup part from a pre-assembly position into an end position. In such case, the stirrup part is rotatably mounted on the first casing, and the second casing has a blocking element which releases pivoting of the stirrup part from the pre-assembly position into the end position only once the two casings are in the pre-assembly position, as described in U.S. Pat. No. 6,540,532 B1."}
{"text": "Microelectronic chips are typically flat bodies with oppositely facing, generally planar front and rear surfaces with edges extending between these surfaces. Chips generally have contacts, sometimes also referred to as pads or bond pads, on the front surface which are electrically connected to the circuits within the chip. Chips are typically packaged by enclosing them with a suitable material to form microelectronic packages having terminals that are electrically connected to the chip contacts. The package may then be connected to test equipment to determine whether the packaged device conforms to a desired performance standard. Once tested, the package may be connected to a larger circuit (e.g., a circuit in an electronic product such as a computer or a cell phone) by connecting the package terminals to matching lands on a printed circuit board (PCB) by a suitable connection method such as soldering.\nMicroelectronic packages may be fabricated at the wafer level; that is, the enclosure, terminations and other features that constitute the package, are fabricated while the chips, or die, are still in a wafer form. After the die have been formed, the wafer is subject to a number of additional process steps to form the package structure on the wafer, and the wafer is then diced to free the individually packaged die. Wafer level processing can be an efficient fabrication method because the footprint of each die package may be made identical, or nearly identical, to the size of the die itself, resulting in very efficient utilization of area on the printed circuit board to which the packaged die is attached.\nA common technique for forming electrically conductive connections between a microelectronic chip and one or more other electronic components is through wirebonding. Conventionally, a wirebonding tool attaches the end of a wire to a pad on a microelectronic chip using thermal and/or ultrasonic energy and then loops the wire to a contact on the other electronic component and forms a second bond thereto using thermal and/or ultrasonic forces."}
{"text": "Field of the Invention\nThe invention concerns a method for evaluating medical image data of an examination subject, as well as an evaluation computer, a medical imaging apparatus and a non-transitory data storage medium for implementing such a method.\nDescription of the Prior Art\nTypically, medical image data are recorded by medical imaging systems and can represent anatomical structures and/or functional processes of the body of an examination subject. An approach for evaluating the medical image data is known in which the medical image data acquired from the examination subject are compared with reference image data of the same examination subject in a healthy state, i.e. in particular in a state prior to the onset of disease symptoms. However, the presence of the reference image data is necessary for this approach, and such data are not available in most cases.\nA further possibility for evaluating the medical image data is known that can be implemented independently of the presence of reference image data. The medical image data acquired from the examination subject are evaluated by semiautomatic or automatic analysis tools. This enables a value of a physiological parameter to be determined on the basis of the medical image data acquired from the examination subject. The thus determined value of the physiological parameter can be compared with a normal value range. In this case the normal value range typically indicates percentiles within which the value of the physiological parameter can be considered normal or non-suspect. A value of the physiological parameter outside of the normal value range, on the other hand, can point to a suspect situation or to the presence of a disease in the examination subject."}
{"text": "In order to quickly review and analyze a video, for example a movie, a recorded sporting event or a news broadcast, a summary of the video can be generated. A number of techniques are known for summarizing uncompressed and compressed videos.\nThe conventional practice is to first segment the video into scenes or ‘shots’, and then to extract low and high level features. The low level features are usually based on syntactic characteristics such as color, motion, and audio components, while the high level features capture semantic information.\nThe features are then classified, and the shots can be further segmented according to the classified features. The segments can be converted to short image sequences, for example, one or two seconds ‘clips’ or ‘still’ frames, and labeled and indexed. Thus, the reviewer can quickly scan the summary to select portions of the video to playback in detail. Obviously, the problem with such summaries is that the playback can only be based on the features and classifications used to generate the summary.\nIn order to further assist the review, the segments can be subjectively rank ordered according to a relative importance. Thus, important events in the video, such as climactic scenes, or goal scoring opportunities can be quickly identified, see, Fujiwara et al. “Abstractive Description of Video Using Summary DS,” Point-illustrated Broadband+Mobile Standard MPEG Textbook, ASCII Corp., p. 177 FIGS. 5-24 Feb. 11, 2003, also “ISO/IEC 15938-5:2002 Information technology—Multimedia content description interface—Part 5: Multimedia Description Schemes,” 2002. After an important video segment has been located, the viewer can use fast-forward or fast-reverse capabilities of the playback device to view segments of interest, see “DVR-7000 Instruction Manual,” Pioneer Co., Ltd., p. 49, 2001.\nAnother technique for summarizing a news video uses motion activity descriptors, see U.S. patent application Ser. No. 09/845,009, titled “Method for Summarizing a Video Using Motion Descriptors,” filed by Divakaran, et al., on Apr. 27, 2001. A technique for generating soccer highlights uses a combination of video and audio features, see U.S. patent application Ser. No. 10/046,790, titled “Summarizing Videos Using Motion Activity Descriptors Correlated with Audio Features,” filed by Cabasson, et al., on Jan. 15, 2002. Audio and video features can also be used to generate highlights for news, soccer, baseball and golf videos, see U.S. patent application Ser. No. 10/374,017, titled “Method and System for Extracting Sports Highlights from Audio Signals,” filed by Xiong, et al., on Feb. 25, 2003. Those techniques extract key segments of notable events from the video, such a scoring opportunity or an introduction to a news story. The original video is thus represented by an abstract that includes the extracted key segments. The key segments can provide entry points into the original content and thus allow flexible and convenient navigation.\nThere are a number of problems with prior art video recording, summarization and playback. First, the summary is based on some preconceived notion of the extracted features, classifications, and importance, instead of those of the viewer. Second, if importance levels are used, the importance levels are usually quantized to a very small number of levels, for example, five or less. More often, only two levels are used, i.e., the interesting segments that are retained, and the rest of the video that is discarded.\nIn particular, the hierarchical description proposed in the MPEG-7 standard is very cumbersome if a fine quantization of the importance is used because the number of levels in the hierarchy becomes very large, which in turn requires management of too many levels.\nThe MPEG-7 description requires editing of the metadata whenever the content is edited. For example, if a segment is cut out of the original content, all the levels affected by the cut need to be modified. That can get cumbersome quickly as the number of editing operations increases.\nThe importance levels are highly subjective, and highly context dependent. That is, the importance levels for sports videos depend on the particular sports genre, and are totally inapplicable to movies and news programs. Further, the viewer has no control over the length of the summary to be generated.\nThe small number of subjective levels used by the prior art techniques make it practically impossible for the viewer to edit and combine several different videos based on the summaries to generate a derivate video that reflects the interests of the viewer.\nTherefore, there is a need to record and reproduce a video in a manner that can be controlled by the viewer. Furthermore, there is a need for specifying importance levels that are content independent, and not subjective. In addition, there is a need to provide more than a small number of discrete importance levels. Lastly, there is a need to enable the viewer to generate a summary of any length, depending on a viewer-selected level of importance."}
{"text": "1. Technical Field\nThe present invention relates to an image processing system, an image processing apparatus control method, and a recording medium storing an image processing apparatus control program.\n2. Background Art\nWith increased computerization of information, image processing apparatuses such as printers and facsimiles used for outputting the computerized information and scanners used for computerizing documents have become indispensable. In most cases, these image processing apparatuses are configured as multifunctional peripherals (MFPs) that can be used as a printer, facsimile, scanner, and copier by implementing an image pickup function, image forming function, communication function, etc.\nOn the other hand, mobile phones have also become highly advanced, and mobile information processing devices such as smart phones and tablet devices that have information processing functions approaching the sophistication of PCs (hereinafter referred to as “mobile devices”) have become popular. These mobile devices generally include a home screen as a starting point for user operations, with applications accessed via icons displayed on the home screen.\nIn these kinds of mobile devices, in order to cooperate with the image processing apparatuses effectively, a technology that displays a scanning configuration screen of the image processing apparatus on the mobile device and instructs the image processing apparatus to perform scanning in accordance with the configuration input on the configuration screen on the mobile device has been proposed (e.g., JP-2008-193265-A)."}
{"text": "The chemical name of levetiracetam, a single enantiomer, is (−)-(S)-α-ethyl-2-oxo-1-pyrrolidine acetamide, its molecular formula is C8H14N2O2 and its molecular weight is 170.21. Levetiracetam is chemically unrelated to existing antiepileptic drugs (AEDs). It has the following structural formula:\n\nSituations in which AED therapeutic drug monitoring (“TDM”) are most likely to be of benefit have been described (Epilepsia, 49(7):1239-1276, 2008). They include (1) when a person has attained the desired clinical outcome, to establish an individual therapeutic concentration which can be used at subsequent times to assess potential causes for a change in drug response; (2) as an aid in the diagnosis of clinical toxicity; (3) to assess compliance, particularly in patients with uncontrolled seizures or breakthrough seizures; (4) to guide dosage adjustment in situations associated with increased pharmacokinetic variability (e.g., children, the elderly, patients with associated diseases, drug formulation changes); (5) when a potentially important pharmacokinetic change is anticipated (e.g., in pregnancy, or when an interacting drug is added or removed); and (6) to guide dose adjustments for AEDs with dose-dependent pharmacokinetics.\nTherapeutic drug management of levetiracetam would serve as an excellent tool to ensure compliance in administering chemotherapy with the actual prescribed dosage and achievement of the effective serum concentration levels. The role of TDM for levetiracetam may also be useful in managing patients that are overdosed.\nTo date no anti-levetiracetam antibody has been produced and no commercially available immunoassay has been developed for levetiracetam."}
{"text": "The implantation, use, and tensioning of orthopedic implanted devices for fusion, stabilization, and fixation of joints, fractures, and other fusions may involve various complicated processes and components. Many of these devices may use invasive exterior tensioning devices, which may provide little control of tension and/or inaccurate means of measuring tension.\nThe present application appreciates that devising orthopedic systems, methods, and apparatuses for fusion, stabilization, or fixation of bones may be a challenging endeavor."}
{"text": "Adenoviruses are double-stranded DNA viruses that have been isolated from a wide variety of avian and mammalian species, including human, swine, cow and sheep. The adenoviruses cause enteric or respiratory infection in humans as well as in domestic and laboratory animals.\nAt least 47 serotypes of human adenoviruses have been described. Reviews of the most common serotypes associated with particular diseases have been published. See for example, Foy H. M. (1989) Adenoviruses In Evans, AS (ed). Viral Infections of Humans. New York, Plenum Publishing, pp 77-89 and Rubin B. A. (1993) Clinical picture and epidemiology of adenovirus infections, Acta Microbiol. Hung 40:303-323. The complete genome sequence of human adenovirus 5 is disclosed in GenBank accession number M73260. The complete genome sequence of human adenovirus 2 is disclosed in GenBank accession number J01917.\nPorcine adenovirus (PAV) infection has been associated with encephalitis, pneumonia, kidney lesions and diarrhea. See Derbyshire (1992) In: Diseases of Swine (ed. Leman et al.), 7th edition, Iowa State University Press, Ames, Iowa. pp. 225-227. It has been shown that PAV is capable of stimulating both humoral response and mucosal antibody responses in the intestine of infected piglets. Tuboly et al. (1993) Res. in Vet. Sci. 54:345-350. Cross-neutralization studies have indicated the existence of at least five serotypes of PAV. See Derbyshire et al. (1975) J. Comp. Pathol. 85:437-443; and Hirahara et al. (1990) Jpn. J. Vet. Sci. 52:407-409. Previous studies of the PAV genome have included the determination of restriction maps for PAV Type 3 (PAV-3) and cloning of restriction fragments representing the complete genome of PAV-3. See Reddy et al. (1993) Intervirology 36:161-168. In addition, restriction maps for PAV-1 and PAV-2 have been determined. See Reddy et al. (1995b) Arch. Virol. 140:195-200.\nNucleotide sequences have been determined for segments of the genome of various PAV serotypes. Sequences of the E3, pVIII and fiber genes of PAV-3 were determined by Reddy et al. (1995) Virus Res. 36:97-106. The E3, pVIII and fiber genes of PAV-1 and PAV-2 were sequenced by Reddy et al. (1996) Virus Res. 43:99-109, while the PAV-4 E3, pVIII and fiber gene sequences were determined by Kleiboeker (1994) Virus Res. 31:17-25. The PAV-4 fiber gene sequence was determined by Kleiboeker (1995) Virus Res. 39:299-309. Inverted terminal repeat (ITR) sequences for all five PAV serotypes (PAV-1 through PAV-5) were determined by Reddy et al. (1995) Virology 212:237-239. The PAV-3 penton sequence was determined by McCoy et al. (1996) Arch. Virol. 141:1367-1375. The nucleotide sequence of the E1 region of PAV-4 was determined by Kleiboeker (1995) Virus Res. 36:259-268. The sequence of the protease (23K) gene of PAV-3 was determined by McCoy et al. (1996) DNA Seq. 6:251-254. The sequence of the PAV-3 hexon gene (and the 14 N-terminal codons of the 23K protease gene) has been deposited in the GenBank database under accession No. U34592. The sequence of the PAV-3 100K gene has been deposited in the GenBank database under accession No. U82628. The sequence of the PAV-3 E4 region has been determined by Reddy et al. (1997) Virus Genes 15:87-90. The transcriptional map and complete DNA sequence of PAV-3 genome was reported (Reddy et al., 1998, Virus Res.58::97-106 and Reddy et al., 1998, Virology 251:414-426). Vrati et al. (1995, Virology, 209:400-408) disclose sequences for ovine adenovirus. U.S. Pat. No. 6,492,343 discloses PAV expression and vaccine systems.\nBAV3 was first isolated in 1965 and is the best characterized of the BAV genotypes, containing a genome of approximately 35 kb (Kurokawa et al (1978) J. Virol. 28:212-218). Reddy et al. (1998, Journal of Virology, 72:1394) disclose nucleotide sequence, genome organization, and transcription map of BAV3. Reddy et al. (1999, Journal of Virology, 73: 9137) disclose a replication-defective BAV3 as an expression vector. BAV3, a representative of subgroup 1 of BAVs (Bartha (1969) Acta Vet. Acad. Sci. Hung. 19:319-321), is a common pathogen of cattle usually resulting in subclinical infection (Darbyshire et al. (1965). J. Comp. Pathol. 75:327-330), though occasionally associated with a more serious respiratory tract infection (Darbyshire et al., 1966 Res. Vet. Sci. 7:81-93; Mattson et al., 1988 J. Vet Res 49:67-69). Like other adenoviruses, BAV3 is a non-enveloped icosahedral particle of 75 nm in diameter (Niiyama et al. (1975) J. Virol. 16:621-633) containing a linear double-stranded DNA molecule. BAV3 can produce tumors when injected into hamsters (Darbyshire, 1966 Nature 211:102) and viral DNA can efficiently effect morphological transformation of mouse, hamster or rat cells in culture (Tsukamoto and Sugino, 1972 J. Virol. 9:465-473; Motoi et al., 1972 Gann 63:415-418). Cross hybridization was observed between BAV3 and human adenovirus type 2 (HAd2) (Hu et al., 1984 J. Virol. 49:604-608) in most regions of the genome including some regions near but not at the left end of the genome. Bovine adenovirus expression and vaccine systems are disclosed in for example, U.S. Pat. Nos. 5,820,868 and 6,001,591.\nFor general background references regarding adenovirus and development of adenoviral vector systems, see Graham et al. (1973) Virology 52:456-467; Takiff et al. (1981) Lancet 11:832-834; Berkner et al. (1983) Nucleic Acid Research 11: 6003-6020; Graham (1984) EMBO J. 3:2917-2922; Bett et al. (1993) J. Virology 67:5911-5921; and Bett et al. (1994) Proc. Natl. Acad. Sci. USA 91:8802-8806.\nFor a review of adenoviruses and adenovirus replication, see Shenk, T. and Horwitz, M. S., Virology, third edition, Fields, B. N. et al., eds., Raven Press Limited, New York (1996), Chapters 67 and 68, respectively.\nAdenoviral vectors are divided into helper-independent and helper-dependent groups based on the region of the adenoviral genome used for the insertion of transgenes. Helper-dependent vectors are usually made by deletion of E1 sequences and substitution of foreign DNA, and are produced in complementing human cell lines that constitutively express E1 proteins. Graham et al. (1977) J. Gen. Virol. 36:59-74; Fallaux et al. (1996) Hum. Gene Ther. 7:215-222; Fallaux et al. (1998) Hum. Gene Ther. 9:1909-1917. An adenovirus E1A region is described in Darbyshire (1966, Nature 211:102) and Whyte et al., 1988, J. Virol. 62:257-265.\nThough E1-deleted viruses do not replicate in cells that do not express E1 proteins, the viruses can express foreign proteins in these cells, provided the genes are placed under the control of a constitutive promoter. Xiang et al. (1996) Virology 219:220-227. Vaccination of animals with adenovirus recombinants containing inserts in the E1 region induced a systemic immune response and provided protection against subsequent challenge. Imler et al (1995) Hum. Gene Ther. 6:711-721; Imler et al. (1996) Gene Therap 3:75-84. This type of expression vector provides a significant safety profile to the vaccine as it eliminates the potential for dissemination of the vector within the vaccine and therefore, the spread of the vector to non-vaccinated contacts or to the general environment. However, the currently used human adenovirus (HAV) based vectors are endemic in most populations, which provides an opportunity for recombination between the helper-dependent viral vectors and wild type viruses.\nUnited States Patent Application Publication 20020019051 discloses chimeric adenoviral vectors. Morsy et al. (1998, P.N.A.S. USA 95:7866-7871); Kochanek et al. (1996 P.N.A.S. USA 93:5731-5736); Clemens et al. (1998, Gene Therapy 3:965-972); and Parks et al.(1996, P.N.A.S. USA 93:13565-13570) disclose adenovirus vectors.\nThere remains a need for improved adenoviral vectors, especially adenoviral vectors for expression of transgenes in mammalian cells, and for the development of effective recombinant adenovirus vectors for use in immunization and expression systems.\nAll references and patent publications disclosed herein are hereby incorporated herein in their entirety by reference."}
{"text": "It has long been known that bacteria play a part in pathogenic dental caries. More recently an abundance of research has implicated bacteria as the causative agents of periodontal disease. Researchers have found a relationship between tooth deposits in periodontal disease. In 1965 a causal relationship was demonstrated between the daily accumulation of dental plaque which causes gingivitis. Others demonstrated in dogs that gingivitis if untreated progresses to periodontitis. Some different organisms are involved in gingivitis than those involved in periodontal disease. However, the organisms are in both cases various strains and types of bacteria.\nTo date, the art of treating caries, gingival and periodontal diseases is primarily surgical although some advances have been made in the control of disease by chemical means. Historically, the use of a toothbrush with dental floss has been recognized as a prevention and a treatment of the initial carious and gingival lesions. However, once the lesions progress beyond their initial stages surgical intervention is the treatment of choice. In the treatment of caries the lesion is excised using a rotating carbide steel or diamond burr mounted in a high speed handpiece. The cavity preparation is designed to accommodate a silver amalgam, gold or composite type of filling material which simulates the physiologic architecture of the tooth. In the treatment of periodontal diseases the goal of treatment is to remove tooth deposits with hand scalers and currettes and to surgically excise the periodontal pocket so the pocket is unavailable to foster growth of microorganisms.\nSignificant advances have been made in the field of preventative cariology through the employment of fluoride containing dentifrices and through the use of viscous polymers which are painted into the pits and fissures of the crown to obliterate these niches where routine oral hygiene practices are ineffective. The purpose of using fluoride as a nutritional supplement or as a topical agent is to incorporate fluoride into the hydroxyappetite crystalline structure of the enamel. This makes for a more symmetrical and perfect crystal structure which is more resistant to acid demineralization.\nPlastic polymeric materials have been shown to be effective in preventing caries by painting them into pits and fissures soon after the eruption of each tooth into the oral cavity. This system works by etching the enamel surface with an acid to leach out surface enamel crystals and make a porous enamel surface. The viscous polymeric materials are painted onto the porous enamel surface and form an intimate mechanical bond with the enamel. Disadvantages of the polymer system include the necessary application of the material immediately after the eruption of the tooth, their propensity to dissolve and abraid and the cost of application.\nThe control of gingival and periodontal diseases via chemical means has had less success; although, many compounds both in literature and the patent have been proposed for inclusion in dentrifices or mouthwashes. At the present time several chemicals are being considered for their efficacy in preventing or treating periodontal diseases, but there is no accepted therapeutic which is universally recommended by the dental profession. Among such compound considered for use are quatenary ammonium salts, ceramide peroxide, chlorhexidine, systemic and topical antibiotics, alexidine and other compounds having bactericidal efficacy. This approach recognizes that killing plaque bacteria controls caries, gingival and periodontal diseases.\nThe bactericidal properties of hydrogen peroxide are due to its dissociation into hydroxyl radicals which are toxic to bacteria. This property led long ago to consideration for its use in an oral dentifrice (U.S. Pat. No. 959,605, U.S. Pat. No. 975,814). However, The peroxide containing dentifrices of the prior art were not universally believed effective bactericides. Recent research had led to the proposal that the reason for their ineffectiveness was the slow dissociation of hydrogen peroxide into free radicals when contained in a typical composition. This dissociation is not rapid enough for a sufficient quantity of bactericidal free radicals to be generated during the application of the composition. An enzymatic method of generating hydrogen peroxide in situ (U.S. Pat. Nos. 4,150,113 and 4,178,362) attempted to overcome this problem by directing the formation of hydrogen peroxide to specific areas. However, this system does not appear to have gained widespread acceptance.\nAn alternate approach to the problem is to disrupt and eliminate the plaque matrix with compounds or enzymes which can affect this result. To this end, a variety of lipases, hydrolases, invertases and proteases, alone and in combination with each other as well as other compounds have been used (U.S. Pat. No. 3,235,460, U.S. Pat. No. 3,194,738, U.S. Pat. No. 4,154,815, U.S. Pat. No. 2,527,686, U.S. Pat. No. 4,155,868). The main drawback to all of these admixtures is the long period of time that they need to be in contact with the tooth surface."}
{"text": "1. Field of the Invention\nThis invention relates to chemical processing and, more particularly, to systems and methods for monitoring attributes of multiple chemical mixtures using a single sensor, and further to systems and methods for controlling attributes of multiple chemical mixtures monitored by a single sensor, the systems and methods being especially configured for use in semiconductor processing.\n2. Description of the Related Art\nThe information described below is not admitted to be prior art by virtue of its inclusion in this Background section.\nThere is much that separates the wide variety of activities that constitute the field of chemical processing. One common thread between chemical processes, however, is the extent to which variations in the attribute values of the chemical mixtures used in these processes can influence process outcomes. Chemical processes are typically designed on the assumption that values of chemical mixture attributes (e.g., temperature, concentration, particle count, and resistivity) can be maintained within predetermined ranges during processing. In many cases, these processes may be further designed on predictions of how chemical mixture attributes will vary within such ranges over time. Generally speaking, the success of a chemical process is dependent on how closely chemical mixture attributes can be maintained to ideal levels, with the degree of closeness necessary for success varying from industry to industry and from process to process. If, however, result-effective chemical mixture attributes are allowed to fall outside of predetermined ranges, a chemical process will often fail to achieve its goals.\nOne field of chemical processing where the importance of the above can be seen is in semiconductor processing. Chemical mixtures are used in a variety of ways in semiconductor processing, and play a particularly important role in etching and cleaning processes. Of the cleaning processes used in semiconductor processing, one of the most prevalent is the standard clean 1 (SC1) cleaning process. SC1 clean is the first step of the traditional RCA clean. SC1 solutions are solutions of water, hydrogen peroxide (H2O2), and ammonium hydroxide (NH4OH) (in order of decreasing typical concentration), and are often heated to between 60 and 85° C. during use. SC1 solutions may additionally include other chemicals, such as chelating agents used to bind up metallic ions present in the solution.\nThe SC1 clean may be used to remove residual organic and metallic contaminants that remain after various processing steps. For example, SC1 solutions are often used to remove residue remaining from chemical-mechanical polishing (CMP) procedures. CMP procedures are considered “dirty” procedures, and as such a great deal of residue often remains on wafer surfaces after such processes are complete. SC1 cleans may be used to remove this residue, reducing the defect probability in the finished product. In a typical post-CMP cleaning sequence, the wafers to be cleaned are immersed in a SC1 solution bath for a specified time to remove residue remaining from the CMP process. Alternately, spray-cleaning methods incorporating SC1 solutions may be used to clean the wafers.\nWhen using SC1 solutions to clean silicon-bearing surfaces (e.g., a single-crystal silicon wafer and any polysilicon deposited thereupon), however, care must be taken to maintain concentrations of the component chemicals within certain values. While a properly balanced SC1 solution will not remove an inordinate amount of the silicon during cleaning, silicon surfaces can suffer chemical attack if concentrations deviate too severely from desired ranges during processing. For example, ammonium hydroxide will, in the absence of hydrogen peroxide, etch silicon. To avoid undesirable attack of silicon surfaces, it is important to maintain hydrogen peroxide concentration values within a SC1 solution at sufficiently high levels during processing.\nUnfortunately, the hydrogen peroxide within an SC1 solution often decomposes over time. Such decomposition can occur, for example, as the result of impurities accumulating within the solution or if the solution temperature rises too far above desired levels. If an SC1 solution is used in which sufficient hydrogen peroxide is not present, the ammonium hydroxide can severely attack silicon surfaces, possibly reducing polysilicon feature sizes beyond acceptable levels or etching the wafer backside in such a way that problems are created further along in the manufacturing sequence. As linewidths reach 0.15 microns and below, it becomes even more important to maintain the component concentrations of an SC1 solution within acceptable ranges.\nTo prevent the occurrence of such undesirable situations, the concentration and other attributes of chemical solutions used in semiconductor processing have historically been checked using various qualification processes. Qualification processes attempt to characterize various parameters of a process, such as how solution properties change with continued use over time. For example, in one type of qualification process a series of test wafers may be sequentially sent through a chemical bath to determine how long the bath can be used in processing before one or more chemical concentrations (or other attributes) of the cleaning solution within the bath fall outside of acceptable ranges. From the data obtained in such a qualification process, the bath life may be estimated, and the overall cleaning sequence can be adjusted accordingly. However, such characterization techniques only estimate concentration from historical data, and as such are not always sensitive to the variations that may occur from run-to-run. In addition, when qualification wafers are being run through a process production wafers are not, so the extensive qualification of a process can reduce the total amount of time that process is available for use on production wafers.\nIn addition, other events often arise during processing for which qualifying processes are inadequate means of prevention. For example, chemical solutions used in semiconductor processing are often mixed from separate sources in a single tank during a process called pour-up. In a typical design, metering pumps are used to supply each constituent chemical of the solution being poured-up. The tanks also may include capacitive sensors to detect solution levels and control the balance of deionized water for pour-ups. In such configurations, the pour-up process may be an automated process that utilizes the metering pumps and capacitive sensors to mix a solution at a desired concentration. If not all the components of the pour-up system are operating optimally, however, bad pour-ups (i.e., situations in which the solution concentrations after pour-up are not at desired levels), can sometimes result. Even worse, one or more of these metering pumps can, from time to time, suffer total failure. Such metering pump failure may result in the chemical supplied from the failed metering pump being absent altogether from the final solution. What's more, metering pump failures are often not immediately detected, and may escape notice until the next qualification process is run. Processing a wafer with such an improperly balanced solution can cause irreversible damage, and if gone undetected, can even result in the loss of several lots of production wafers.\nTo avoid such problems, some manufacturers have resorted to placing concentration sensors on each tank supplying process chemicals to ensure that solution concentrations are within acceptable levels before the solutions are used in the processing of production wafers. Unfortunately, the concentration sensors required to reach desired accuracy levels are expensive, often costing $50,000 dollars or more. Furthermore, such sensors typically need to be rigorously maintained, further increasing their associated cost of ownership. In many cases, the cost of installing and maintaining concentration sensors on every chemical mixture that would profit from monitored outweighs the benefits obtained from concentration monitoring. Consequently, concentration monitoring may only be implemented in select tanks or forgone entirely.\nTherefore, it would be desirable to design a monitoring system capable of monitoring an attribute (e.g., concentration) of multiple chemical mixtures within multiple chemical vessels during processing that did not require a separate concentration sensor for each chemical vessel being monitored."}
{"text": "This invention relates to laminates incorporating a fibrous layer, and particularly to disposable absorbent products such as panty shields, sanitary napkins, diapers, absorbent towels, and bed liners.\nPanty shields consist essentially of three layers: a permeable layer through which liquid passes, a layer of fibers for absorbing the liquid, and a moisture-barrier layer beneath the fibers. Some means must be provided for holding the fibers together to form a layer so that the laminate does not fall apart, and the layer of fibers must be bonded to the adjoining layers.\nSeveral approaches have been taken in the prior art. One common technique is to seal the fibers within an envelope formed by sealing the permeable and barrier layers together at their perimeters. These layers are typically thermoplastic, and are bonded to each other by application of heat (or with adhesive).\nAnother approach has been to add thermoplastic powders, fibers, or fibrils to the fibrous layer and to apply heat during lamination to form a supporting web of thermoplastic material throughout the fibrous layer as well as a bond between the thermoplastic web and adjoining, thermoplastic permeable and barrier layers. Typically, the laminate is embossed during application of heat so that bonding occurs primarily in embossed areas.\nCrosslinkable latex adhesives have also been employed. The layer of fibers is impregnated with such adhesives during lamination, and the laminate is then cured in drying ovens to form bonds between fibers as well as between the layer of fibers and the adjoining permeable and barrier layers. This is commonly done stepwise in separate operations. Embossing has also been used in such constructions to cause the latex adhesive to migrate into the embossed areas prior to curing.\nAbsorbent fibrous webs have also been made on specialized equipment that consolidates loose fibers using such techniques as spunlacing, mercerizing, or needle punching to mechanically entangle the fibers into a web. Uniform fibers are generally required, and the mechanical entanglement step must be done off-line (i.e., in a separate manufacturing process not part of the process for assembling the panty shield or other disposable absorbent article).\nHot melt adhesives have been used in the manufacture of laminates. A known technique for applying such adhesives is to apply a coating of the adhesive to one or more layers prior to forming the laminate."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a gas lift plunger assembly for use in a gas-producing well and, more particularly, is concerned with pad plunger assembly with one-piece locking end members.\n2. Description of the Prior Art\nGas-producing wells typically employ a plunger disposed within tubing of the well and capable of traveling vertically in the tubing as the well is cycled between shut-in and opened conditions in a manner well-known to one of ordinary skill in the art. The plunger is freely movable vertically in the well tubing and is adapted to rise vertically under the force of sufficient gas pressure to drive or lift the plunger and a slug of liquid, such as oil, above it to the surface while isolating the base of the liquid slug from the gas which lifts the plunger. The plunger falls by gravity back down the tubing of the well after the slug of liquid has been delivered to the surface and the gas pressure in the well tubing has decreased due to the transmission of gas from the well tubing to a suitable storage location.\nOne prior art plunger assembly, called a pad plunger, includes an elongated rigid non-flexible central rod or mandrel and a pair of end members attached to the opposite ends of the mandrel. The end members of the plunger have larger diameters than the mandrel and are internally threaded to permit their removable attachment to the externally threaded opposite ends of the mandrel.\nFrom the foregoing description, it is readily apparent that the pad plunger is made up of an assembly of various parts. It is very important to lock the end members on the plunger so that the pad plunger does not come apart while in the well. The cost of retrieving the plunger parts from the well can range from $1000 to $5000 not including lost production.\nDrawbacks exist, however, with the above-described prior art plunger in the manner used to lock the end member on the mandrel. Locking subassemblies 10, such as the one shown in FIG. 1, that have been used are very labor intensive and have required numerous parts, namely, a locking ball 12 and one or more, usually three, locking rings 14 per plunger. The plunger mandrel 16 and end member 18 have to be drilled to provide a hole 20 extending through the side of the end member 18 and into the mandrel 16 so that the locking ball 12, such as xe2x85x9c inch in diameter, can be dropped in the hole 20 and then the locking ring 14 applied about an annular recess 22 in the end member 18 so as to overlie the hole 20 and retain the ball 12 in the hole. The presence of the ball 12 in the hole 20 prevents the end member 18 from turning and unthreading from the threaded end of the mandrel 16. However, corrosion problems typically occur behind the rings. Over a short period of time the rings become loose causing premature plunger wear. Therefore, the locking rings are being welded in addition to being crimped, adding more time for labor.\nConsequently, a need still exists for an innovation which will provide a solution to the aforementioned problems in the prior art without introducing any new problems in place thereof.\nThe present invention provides a pad plunger assembly with one-piece locking end members designed to satisfy the aforementioned need. The solution of the present invention to the aforementioned problems is to provide one-piece locking end members with features which require less labor usage and eliminate multiple parts.\nAccordingly, the present invention is directed to a pad plunger assembly which comprises: (a) an elongated central mandrel having a pair of opposite ends with fastening elements thereon; (b) a pair of annular collars each attached to and extending about the mandrel inwardly adjacent to one of the opposite ends thereof, each of the collars having formed therein a plurality of recesses circumferentially spaced from one another about the collar; (c) a pair of end members each having an inner end with an interior bore and fastening elements therein and being removably fastened to one of the opposite ends of the mandrel; and (d) a pair of inner annular flanges each attached to and extending axially from the inner end of one of the end members and adapted to be received over one of the collars such that, for locking the end members onto the mandrel, the end members are fastened onto the opposite ends of the mandrel and then the annular flanges are crimped so as to force circumferentially displaced portions thereof into the recesses of the collars thereby locking the end members in place on the opposite ends of the mandrel such that the end members cannot then be unfastened from the mandrel nor become loosened over time.\nMore particularly, the mandrel is an elongated rigid non-flexible solid rod. The collars are integrally connected to the mandrel. Each of the inner annular flanges is integrally attached to the inner end of one of the end members. The fastening elements on the ends of the mandrel are external threads. The fastening elements in the interior bores of the end members are internal threads which threadably receive the external threads on the ends of the mandrel when the end members are fastened thereon. The mandrel, collars, end members and flanges are made of a metal.\nAlso, the assembly includes at least one subassembly of spring-loaded interlocking pad sections disposed about the mandrel and extending over the mandrel between the collars.\nThese and other features and advantages of the present invention will become apparent to those skilled in the art upon a reading of the following detailed description when taken in conjunction with the drawings wherein there is shown and described an illustrative embodiment of the invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a coordinates input apparatus and, more particularly, to a coordinates input apparatus which detects coordinates of an indication point from a vibration propagation time on a vibration propagating plate and has a structure in which vibration proof materials are arranged in edge portions of the vibration propagating plate.\n2. Related Background Art\nHitherto, there has been known a coordinates input apparatus in which a vibration is input to a vibration propagating plate by a vibration pen having therein a piezoelectric transducing element or the like, the input vibration is detected by a plurality of sensors attached to the vibration propagating plate, the vibration propagation times are measured, and the coordinates of the input point are detected.\nIn such a coordinates input apparatus, there is a construction in which the edge portions of the vibration propagating plate are supported by vibration proof materials so as to prevent a situation such that the input vibration is reflected by the edge portions of the vibration propagating plate and an error is caused in the detection by the vibration sensor due to the reflected waves.\nThe conventional vibration proof material has mainly been developed to prevent noises and is largely classified as a material for use as a countermeasure for the air sound and a material for use as a countermeasure for the solid sound. Therefore, in the case of using the conventional vibration proof material in the above application fields, the material as a countermeasure for the solid sound is used. As a conventional damping material as a countermeasure for the solid sound, a damping sheet for a thin plate as shown in FIG. 8A, a paint, or the like has been known.\nIn FIG. 8A, reference numeral 8' denotes a vibrating plate whose vibration should be damped and 7 indicates a damping sheet. The vibrating plate 8' is made of a metal plate such as thin steel plate, aluminum plate, or the like, a resin plate, a glass plate, or the like.\nIn such a structure, by adhering the damping sheet 7 onto the vibrating plate 8' which is vibrating, the vibration of the plate 8' is decreased by using a vibration attenuation of the damping sheet 7, so that the noises can be reduced.\nAs materials for commercially available damping sheets 7, there are a polyvinyl chloride resin, atactic polypropylene, polyethylene vinyl acetylate, a styrene-butadiene rubber, a silicon rubber, a cement paste, and the like. Further, the products which are obtained by adding or mixing a plasticizer, a stabilizer, a softener, metal powder of lead, iron, or the like, quarts sand, asphalt, and the like to the above materials are used. As a molding shape of those materials, there is a sponge-like porous shape.\nThe conventional damping sheets as mentioned above are considered to suppress the vibration of the whole plate mainly by adhering the damping sheet onto the whole plate which vibrates.\nTherefore, even if the conventional damping sheet is attached to the periphery of the vibration propagating plate 8 as shown in FIG. 8B in order to suppress the reflected waves at the edge portions of the vibration propagating plate, the reflected waves cannot be sufficiently reduced.\nFIG. 8B is a diagram showing the vibration propagation in the case where the damping sheet has been attached to the peripheral portion of the vibration propagating plate. FIG. 8C is a diagram of a part of a cross sectional view of FIG. 8B and shows a state of the reflected waves.\nIn the diagrams, reference numeral 8 denotes a vibration propagating plate; 7 the damping sheet; 3 a vibration pen as a source to apply a vibration; a wave showing a vibration which is propagated from a vibration applied point as a contact point between the vibration pen 3 and the vibration propagating plate 8; B a reflected wave at the edge surface of the vibration propagating plate; and C a reflected wave at the boundary surface in the portion where the damping sheet has been attached.\nThe damping effect of the conventional damping sheet 7 is large for the vibration (also including a natural vibration, i.e., resonant vibration) which occurs in the whole plate as mentioned above. However, as shown in FIG. 8B, an adequate damping effect is not derived for a vibration which has been applied to a region where no damping sheet is attached and which propagates as a progressing wave from such a region. Therefore, as shown in FIG. 8C, although the vibration is slightly attenuated in a portion D to which the damping sheet 7 has been attached, the reflected wave B at the edge surface of the vibration propagating plate cannot be sufficiently suppressed.\nFurther, by attaching the conventional damping sheet 7, the reflected wave C is newly generated at the interface surface where the damping sheet has been attached. Consequently, if the conventional damping sheet has merely been attached to the peripheral portion of the vibration propagating plate of the coordinates input apparatus using an elastic wave, the foregoing two reflected waves are generated and become noises when a direct wave from the vibration applying source is detected, so that a detecting precision is deteriorated.\nWhen considering the conventional example of the vibration proof material (damping sheet) attached to the vibration propagating plate mentioned above, it is necessary to pay an attention to the following two points with respect to the positional relation between the vibration sensor (hereinafter, simply referred to as a sensor) and the vibration proof material which is attached to the peripheral portion of the vibration propagating plate in the coordinates input apparatus in which a vibration which has been input from the vibration pen is detected by a plurality of vibration sensors attached to the vibration propagating plate and the position of the vibration pen is detected from the vibration propagation times until the vibration sensors. One point relates to an influence which is exerted on the vibration propagation by the attenuation of the vibration proof material. Another point relates to an influence by the reflected wave which is generated by the surface on which the vibration proof material has been attached. The same applicant as the present invention has hitherto proposed the following apparatuses with regard to the above two points. In the \"Coordinates input apparatus\" disclosed in Japanese Patent Application No. 62-67856, by attaching sensors a predetermined distance away from a vibration proof material (two wavelengths) or longer, the influence exerted on the vibration propagation by the attenuation of the vibration proof material is reduced. In the \"Coordinates input apparatus\" disclosed in Japanese Patent Application No. 61-251598, by attaching sensors so as to be away from a vibration proof material by a predetermined distance (reflection interference distance) or longer, the influence by the reflected wave which is generated by the vibration proof material attaching surface is reduced. Further, in the \"Coordinates input apparatus\" disclosed in Japanese Patent Application No. 61-251599, by attaching sensors onto the vibration proof material attaching boundary surface, the influence of the reflected wave which is generated by the vibration proof material attaching surface is reduced.\nIn the above conventional apparatuses, however, there are drawbacks such that in the case of attaching the sensors so as to be away from the vibration proof material by a predetermined distance or longer, outer dimensions of the whole apparatus are too large, while in the case of attaching the sensors onto the vibration proof material attaching boundary surface, the vibration proof material and the sensors cannot be attached onto the same surface of the vibration propagating plate and the thickness of apparatus increases and the like."}
{"text": "Heat-labile enterotoxin (LT), produced by enterotoxigenic strains of E. coli, and cholera toxin (CT), produced by V. cholerae strains, are the causative agents of traveller's diarrhoea and cholera, respectively [Spangler (1992) Microbiol Rev 56:622]. LT and CT show 80% homology in the primary structure and an identical tertiary structure. They are composed of two functionally distinct domains: the enzymatically-active A subunit and the B pentamer, which contains the receptor-binding site. The A subunit ADP-ribosylates the target protein Gs, a GTP-binding protein which regulates the intracellular levels of cAMP [Rappuoli & Pizza (1991), in Sourcebook of bacterial protein toxins, Academic Press NY]. Enhancement in cAMP levels can alter ion transport, inducing secretion of water and chloride ions into the intestine.\nCT and LT are both powerful immunogens and potent mucosal adjuvants when co-administered with antigens at the mucosal level [eg. Jackson et al. (1993) Infect Immun 61:4272; WO95/17211]. Immunogenicity and adjuvanticity of wild-type CT and LT have been extensively studied in animals [eg. Rollwagen et al. (1993) Vaccine 11:1316], but their toxicity has precluded their use in humans. In an attempt to overcome the problems generated by the use of active holotoxins, two different approaches have been followed, one based on the use of the B pentamer, the non-toxic domain of the holotoxin [eg. Holmgren et al. (1992) Vaccine10:911], and the other based on the generation of genetically detoxified derivatives of LT and CT [eg. Fontana et al. (1995) Infect Immun 63:2356]. Site-directed mutagenesis on both A and B subunits has provided a tool to explore the basis of the immunological and adjuvant responses induced by these molecules.\nExamples of such experiments can be found in:                Harford et al. [Eur J Biochem (1989) 183:311] made LT-A carrying Ser-61-Phe and Gly-79-Lys substitutions.        Tsuji et al. [J Biol Chem (1990) 265:22520] produced LT-A carrying a Glu-112-Lys substitution.        Burnette et al. [Infect Immun (1991) 59:4266] produced CT-A carrying a Arg-7-Lys substitution. This work can also be seen in WO92/19265.        WO93/13202 described non-toxic CT and LT carrying mutations at Val-53, Ser-63, Val-97, Tyr-104, and Pro-106        \nA mutant in the receptor binding site of the B subunit of LT, the G33D mutant, has been reported to lack the immunological properties of the native B subunit [Nashar et al. (1996) PNAS 93:226], suggesting that binding to the receptor is important for the immunogenicity. It has also been shown that non-toxic derivatives of LT carrying mutant A subunits retain the immunological properties of wild-type LT [eg. Magagnoli et al. (1996) 64:5434], suggesting that ADP-ribosylation activity is not essential for immunogenicity.\nIn relation to adjuvanticity, much data has been generated but many questions remain unanswered. Some studies have suggested that LT-B and CT-B have adjuvant activity, but the conclusions drawn have been compromised by contamination with active toxin [Wilson et al. (1993) Vaccine 11:113]. Studies with recombinant LT-B and CT-B have suggested that they behave as poor mucosal adjuvants [eg. Douce et al. (1997) Infect. Immun. 65:2821]\nAttempts to define the role of ADP-ribosylation activity in the adjuvanticity of LT has generated conflicting results. Lycke et al. [Eur J Immunol (1992) 22:2277] have described a non-toxic LT derivative (LT-E112K) which, when administered with KLH by oral route in mice, lacked the adjuvant properties of wild-type LT, thus suggesting that adjuvant activity is linked to ADP-ribosylation activity. LT derivatives such as LT-K7 and LT-K63 [eg. Douce et al. (1997) supra; Douce et al. (1996) PNAS 92:1644; DiTommaso et al. (1996) Infect Immun 64:974], however, which are devoid of toxicity both in vitro and in vivo, have been shown to be able to elicit an antibody response against a co-administered antigen in intranasally immunised mice. LT-K63 has been shown to induce measle-specific CTL response after intranasal immunisation with a synthetic peptide [Partidos et al. (1996) Immunol 89:483], and strongly enhances protection against H. pylori following intragastric immunisation with H. pylori antigens. The antibody titres induced by these non-toxic LT mutants were lower than those obtained with wild-type toxin, however, and were only detected after two mucosal immunisations [Douce et al. (1997) supra].\nIt is an object of the invention to provide forms of LT which are detoxified, so that they might be suitable for use in humans, but which retain the adjuvant and immunogenic properties of LT as far as possible."}
{"text": "UMTS (Universal Mobile Telecommunications System) is a more general term for the 3G, (third generation) telecommunications system based on the WCDMA high capacity radio interface. GSM is the most widespread 2G (second generation) telecommunications system based on TDMA (Time Division Multiple Access) radio. The goal of the GPRS system is to provide global layer 2 connectivity from a cellular mobile terminal (MT, sometimes also referred to as mobile station (MS) or user equipment (UE)) using 2G or 3G radio technology (e.g. GSM, American TDMA, UMTS, GERAN (GSM/EDGE Radio Access Network) to an external packet data network. GPRS can support various layer 3 protocols (e.g. IPv4; IPv6; PPP (Point-to-Point Protocol)).\nThe main nodes of a GPRS network are SGSN (Serving GPRS Support Node) and GGSN (Gateway GPRS Support Node). SGSNs are the nodes serving the MT. Each SGSN supports GPRS for GSM and/or UMTS. GGSNs are the node handling the interworking with PDNs (Packet Data Networks). Signaling and data are exchanged between SGSN and GGSN or SGSN and SGSN using the GTP (GPRS Tunneling Protocol) protocol. GTP protocol handles mobility, and creation, modification and deletion of GTP tunnels, as well as transfer of user data between GSNs. GTP allows multi-protocol packets to be tunneled between GSNs and between an SGSN and the UMTS Terrestrial Radio Access Network (UTRAN, not shown) through which a connection to the concerned MT is established. Other systems components need not be aware of the GTP. Typically, two IP addresses are used for a single tunnel, one for the GTP control message (i.e. signalling) and one for the GTP user packet (i.e. carrying user data).\nAt GPRS attach, the SGSN establishes a mobility management (MM) context containing information pertaining to e.g. mobility and security for the concerned MT. At PDP context activation, the SGSN establishes a PDP context, to be used for routing purposes, with the GGSN the subscriber will be using A GPRS attached mobile terminal (MT) can be assigned either a static or dynamic IP address (referred to also as PDP address). The static address is assigned by the Home Public Land Mobile Network (HPLMN) operator at the time of subscription. The dynamic IP address can be allocated by a GGSN (Gateway GPRS Support Node) of either the HPLMN or the visited PLMN (VPLMN) operator at PDP (Packet Date Protocol) context activation time. In addition to address allocation, the GGSN implements the forwarding of IP packets from a GTP (GPRS Tunneling Protocol) tunnel to a packet data network (PDN) and vice versa. There are two kinds of PLMN backbone networks, an Intra-PLMN backbone network and an Inter-PLMN backbone network. The Intra-PLMN backbone network is a private IP network intended for packet domain data and signaling within a PLMN only, while the Inter-PLMN backbone is used for roaming from one PLMN to another (via Border Gateways). Serving GPRS Support Nodes (SGSNs) and GGSNs use the Intra-PLMN backbone to exchange GPRS domain data and signalling.\nDuring roaming, both the Intra-PLMN backbone of the home and visited networks are used in addition to the Inter-PLMN backbone. When a subscriber is roaming to another PLMN, i.e. a VPLMN, the user needs to first attach to the network. In GPRS attach, the MT informs the Serving SGSN of its intention to connect to the network by giving information about its identity, capability and location. The SGSN then checks the MT's identity and performs an authentication procedure in order to secure the transmission path. The attachment is completed after the SGSN has received roaming subscriber data from the Home Location Register (HLR) of the subscriber's HPLMN and finished a location update procedure. After the GPRS attach, the MT sends an ‘Activate PDP context’ request, in which the Access Point Name (APN) is a reference to the GGSN Access Point (AP) to be used in either the home or visited PLMN. The SGSN selects the GGSN based on a PDP context subscription record and sends the context data to a selected GGSN. The GGSN then routes the packets to the appropriate Packet Data Networks (PDN).\nWhen a subscriber is roaming in the VPLMN, there are the following two possibilities for GGSN selection. Firstly, the home network GGSN can be used via the Inter-PLMN backbone and BGs. The home GGSN then routes the packets to their destination. Secondly, a visited domain GGSN can be used for routing the packets from the VPLMN to their destination directly through a packet data network (PDN), such as the public Internet.\nIt should be noted that there are two levels of IP addressing:                the user IP address corresponding to the packets carried over GTP protocol. The corresponding IP address is referred to as PDP address or user address; and        the network IP address corresponding to the packets carried below GTP protocol. The corresponding IP addresses are the node IP addresses used to exchange GTP packets between GSNs. These IP addresses might also be used for network operation such as charging or O&M.        \nUser and network addresses are independent of each other thanks to the GTP protocol, and could both be either IPv4 or IPv6.\nThe GPRS backbone nodes of the second (2G) and third (3G) generation may optionally use an IPv6 based addressing for network addresses. However, existing specifications allowing the use of IPv6 addresses, do not define how to maintain backward compatibility with IPv4 based nodes. An SGSN should know in advance that the GGSN selected supports IPv6 addresses before inserting an IPv6 address e.g. in a create PDP context request message.\nFurthermore, another problem arises with the known procedures. If a MT moves from an IPv6 capable SGSN connected to an IPv6 capable GGSN to an IPv4 only SGSN, the communication will get lost as the new SGSN cannot use the IPv6 address transferred. Such a scenario is very realistic in particular when two operators with equipment from different manufacturers (or just different software release) have roaming agreements (e.g. national roaming). It should be noted that existing IPv4-to-IPv6 transition mechanisms do not apply here as they do not affect IP addresses carried in GTP.\nIn addition, a practical requirement is that the protocol changes have to be done so that nodes based on older version of specifications (and so not supporting enhancement proposed here), have to continue interworking with new nodes supporting the proposed enhancement proposed here."}
{"text": "The invention relates to controls for automobile carwash installations. More particularly, it relates to controls which activate a plurality of subsystems disposed along a conveyor bearing automobiles to be washed through timer-activated logical controls.\nThe art of constructing and using carwash installations in which an automobile is carried along a conveyor past work stations performing specific tasks--whose totality adds up to a complete cleaning job on the vehicle--are well known in the prior art. Such conventional car washes are based on sensors associated with each, or most, task-performing devices which sense the presence, or approach, of a vehicle and activate the appropriate function.\nSuch sensors are commonly wand-operated limit switches, proximity switches, or guide rollers attached to particular components, such as rotary brushes. The use of such sensors in the inhospitable environment produced by the water, steam, soap and chemical agent saturated confines of the carwash leads to a large number of maintenance problems and can result in the inefficient use of scarce resources, such as hot water, to overcome positioning errors and other factors associated with mechanical and electro-mechanical direct-acting controls.\nEven greater difficulties are encountered in carwash control systems of conventional construction when it is desired to selectively activate certain work stations for particular vehicles; for example by dispensing hot wax where the driver is willing to pay for the higher charges associated with such service.\nWhile proximity sensors and photoelectric devices avoid some of the problems associated with mechanically actuated position indicators, such as possible marring of automobile surface finishes or interfering with external trim, the large variety of shapes, sizes and constructional materials in automobiles make them less accurate and reliable in providing accurate position signals. For these reasons, the use of sensors along the carwash conveyor as parts of the operational control system are attended by many practical difficulties and problems.\nIt is, therefore, a primary object to the invention to provide controls for a carwash installation incorporating a plurality of active devices disposed alongside a conveyor, in which the position of the individual vehicles is tracked by logical devices, rather than determined by position sensors.\nIt is a further object of the invention to provide controls for a carwash installation capable of performing a wide range of optional tasks on individual automobiles as they travel along the conveyor, through the assignment of a dedicated program control unit to each vehicle as it enters the installation.\nIt is also an object of the invention to teach the construction of program control units, and associated control equipment, resulting in economical and reliable control systems for a carwash environment."}
{"text": "The present invention relates generally to seats and more specifically to a hydro-formed seat adjuster system.\nIt is well known in the automotive vehicle seating industry to employ seat adjusters which carry a seat support frame upon which are mounted a seat bottom cushion and a seat back cushion. The seat adjusters can be manually operated or automatically driven by electric motor actuators. Many traditional power seat actuators employ a front torsion tube and a rear torsion tube. Motion links are typically fusion welded onto each torsion tube. These torsion tubes control lifting and tilting of the seat adjuster, seat frame and cushions. One such conventional assembly is disclosed within U.S. Pat. No. 5,899,428 entitled “Vehicle Power Seat Adjuster with Self-Locking Track Assemblies” which issued to Gauger on May 4, 1999, which is incorporated by reference herein."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device and a method of manufacturing the same.\n2. Description of the Related Art\nA micro-electro-mechanical system (MEMS) display is a display expected to replace a liquid crystal display (see Japanese Patent Application Laid-open No. 2008-197668). This display differs from a liquid crystal shutter type display utilizing polarization, and displays an image by opening and closing a light transmissive window using a mechanical shutter system. A shutter is formed of a thin film. Vertical and horizontal sizes of one shutter forming one pixel are in the order of several hundred micrometers, and a thickness thereof is in the order of several micrometers. One shutter is opened/closed to enable ON/OFF operation for one pixel. The shutter is operated by an electrostatic attractive force.\nThe shutter is arranged in a space surrounded by a sealing member between a pair of light transmissive substrates, and the space is filled with oil. The oil is used to prevent a spring for driving the shutter from sticking, and to reduce a difference in optical refraction index with respect to the light transmissive substrates.\nThe oil is injected from an injection port corresponding to an opening of the sealing member, and the injection port is encapsulated by a resin after the oil injection. For high-speed open/close operation of the shutter, oil having low viscosity is desired, but in this case, the speed of oil when passing through the injection port increases to cause a problem of damaging the shutter."}
{"text": "1. Field of the Invention\nThe invention relates to the production of silicon bodies and somewhat more particularly to a process and apparatus for producing plate-shaped silicon bodies useful in solar cell manufacture.\n2. Prior Art\nIn the production of solar cells comprised of large-surface silicon slabs or plates, the cheapest available silicon material is desired. The quality of the silicon material which is used in producing solar cells, especially the crystal quality thereof, is not overly critical, for example, as in instances where silicon material is used to manufacture electrical components.\nA typical prior art process of producing silicon slabs for solar cells comprises growing a large cylindrical monocrystal of silicon by slowly withdrawing a seed crystal from a silicon melt and slicing this monocrystalline rod into thin slabs or the like. The flat surfaces of each slab is then polished or otherwise treated to remove surface irregularities.\nGerman Pat. No. 1,555,916 suggests that plate-like silicon bodies may be produced by heating a silicon rod sufficiently to render such rod ductile and subjecting such ductile rod to a rolling or forging process so as to convert the rod into a flat sheet-like body from which silicon plates or the like may be severed. In this process, the starting rod is a monocrystal so that the ultimately produced silicon plates retain at least some monocrystalline characteristics during the transformation of the rod into a ductile condition and during the subsequent rolling or forging. While substantial crystal disorder in the so-processed structures occurs, such disorder is readily removed by subsequent annealing processes."}
{"text": "1. Field of the Invention\nThe present invention relates to a temperature compensating device and satellite signal receiving system, and more particularly, to a temperature compensating device and satellite signal receiving system which compensates for effects of ambience temperature variation through current stabilizing, to ensure normal system operation.\n2. Description of the Prior Art\nTemperature is a critical factor in characteristics of electronic component; therefore, an ideal electronic product must be able to maintain a stable performance within a reasonable temperature range. To meet this requirement, electronic systems have temperature compensating mechanisms to compensate for effects of temperature changes.\nIn a satellite receiver, for example, after an antenna receives radio frequency signals transmitted by a satellite, a radio frequency amplifier first amplifies the received radio frequency signals, and then a mixer utilizes an oscillating signal of a local oscillator to lower the radio frequency signals to an intermediate or baseband frequency. In this process, abnormal operations in the mixer or the local oscillator would cause demodulation or decoding failure of analog or digital circuits in subsequent stages, leading to malfunction of the satellite receiver. A key influencing factor in operations of the mixer or the local oscillator is ambience temperature change, mainly because satellite receives are often situated outdoors and unsheltered; and in regions where outdoor temperature differences may exceed 50 degrees Celsius, such vast temperature ranges may cause changes in characteristics of electronic components (e.g. resistance values, capacitance values, conduction impedance of transistors, etc.), leading to malfunction of the mixer or the local oscillator.\nSpecifically, the mixer principally operates by utilizing a local oscillator to multiply the radio frequency signals and the local oscillating signal, and then filtering out unneeded signals with a band-pass filter. The local oscillator has a basic structure of a loop constructed from band-pass filters and amplifiers, mainly for repeatedly amplifying noise signals to acquire needed oscillations, and the band-pass filter filters signals at designed frequencies. As is well known by the industry, a multitude of electronic components, e.g. transistors, capacitors, resistors, etc, are required to implement the mixer and the local oscillator, thus making them prone to temperature differences.\nConventional temperature compensating mechanisms are usually implemented through fixed voltages, which have limited stabilizing effects on transistors, thus a fixed current temperature compensating mechanism is required."}
{"text": "Modern communication service providers frequently offer a wide variety of communications services to users. For example, users often purchase, subscribe to, or otherwise obtain voice calling, text messaging, or high speed data access services from communication service providers. In addition, a service provider may claim a certain level of data transfer speed when marketing or branding a high speed data access service.\nA service provider making such a speed claim with regard to an offered high speed data service may desire verification of the actual experienced data rates of a user. While the service provider may monitor certain network conditions from the perspective of the communication network, such measurements merely provide a rough estimate of actual user experience. Unfortunately, existing methods to determine user experience, such as drive testing the network or providing customer satisfaction surveys, are not exhaustive enough to capture a significant data set that accurately represents the entire network."}
{"text": "The development of new methods for preparing three-dimensional structures of the nano- or the micrometer scale is of great interest in various technological fields such as biotechnology, microelectronics, microfluidics and biomaterials.\nIn microelectronics, research have focused on the development of new approaches for miniaturizing classical 2D integrated circuits into smaller 3D-structures in which interconnected silicon wafers (3D-IC (3D-integrated circuits)) or chips (3D-packaging) are stacked. The so-called “3D-technology” enables to provide greater functionality and higher component density while decreasing the form factor. One of the biggest challenges of 3D-technologies remains the creation of vertical electrical connections between stacked chips or stacked silicon wafers using nano- or microstructures such as nanowires. Nowadays, microelectronic devices are mainly available by conventional top-down processes such as photolithography. But, for the foreseeable future, such technologies are believed to remain inappropriate for the large-scale production of microelectronic devices.\nNanotechnology research thus put a great emphasis on the “Bottom-Up” strategies. These strategies are based on the self-assembly of molecular building blocks for creating three dimensional nano- or microstructures to be integrated in larger and functional microdevices. Appropriate molecular building blocks encompass inorganic compounds able to form crystals as well as biological molecules such as DNA, peptides and proteins. In that respect, proteins are very attractive building blocks because of their physical sizes, their highly specific interactions, their high degree of organization and their ability to be chemically functionalized or coated.\nSeveral publications describe nanowires obtained from protein fibres.\nScheibel et al. (PNAS, 2003, 1000, 4527-4532) describe the use of self-assembling amyloid fibres from Saccharomyces cerevisiae prion determinant to construct nanowire elements. The said nanowire elements were subsequently metallized with silver or gold so as to produce conductive metal wires of 100 nm wide.\nPatolsky et al. (Nature Materials, 2004, 4, 692-695) describe gold wires (1-4 μm long and 80-200 nm high) exhibiting high electrical conductivity which are obtained by polymerization of G-actin labelled with Au (gold) nanoparticles followed by the catalytic enlargement by metallization process. In the same way, Yao et al. (Analytical and Bioanalytical Chemistry, 395, 1563-1566) describe semiconductor nanowires obtained by conjugation of quantum dots to actin filaments.\nHowever, few studies relate to methods for actually preparing three-dimensional protein structures more complex than nanowires.\nNakamura et al. (MEMS proceedings, 2007, Japan) describe that, when G-actin and fascin solution was confined and polymerized in Polydimethylsiloxane (PDMS) micro chambers, the shape of the actin bundles may follow the geometry of the chambers. It clearly appears that such a method is quite limited and only enables to obtain very few flat actin structures.\nBrough et al. (Soft Mater, 2007, 3, 541-546) has proposed a method for preparing three-dimensional actin structures with nano-resolution patterns in the XY direction and on the micron pattern on Z direction while controlling the precise location of said structures on substrates such as silicon wafers. This method is based on the use of streptavidin nanopatterns for immobilizing biotinylated gelsolin-F-actin complex and subsequently promoting the polymerization of actin filaments. Gelsolin is indeed a capping agent which binds the barbed of actin filament. Consequently, the polymerization of actin filaments only occurred from pointed ends and is thus very slow. The resulting polymerized actin filaments are unbranched and display a length of at most 3-4 μm. Even if the actin filaments grew in out-of-plane direction from the nanopatterns, no well-orientated and well-defined three-dimensional actin structure could be actually observed.\nHuang et al. (Langmuir 2006, 22, 8635-8638) describes a method similar to that of Brough et al. This method is based on the immobilization of pre-formed actin filaments on patterns. The immobilized actin filaments are then aligned along a desired direction using an electric field.\nUS 2005/0106629 reports a method for preparing a three-dimensional actin column based on the use of a first surface (nucleating plate) patterned with nucleating agent and a second surface (binding plate) patterned with capture agent. As explained in Example 3 of this document, the two patterned surfaces are brought into contact and aligned within a small fluid chamber under an inverted microscope. The fluid chamber is filled with a polymerization mixture to initiate column growth from nucleating plate. For enabling the growth of the actin column, the binding plate is moved away from the nucleating plate at a rate approximately equal to the rate of the column growth but not faster. In other words, the column grows as a bundle of parallel actin filaments, each actin filament being in contact by its first end to a nucleating agent on the nucleating plate and being connected by its second end to a capture agent on the binding plate. Even if said method virtually enables to control the location of the three-dimensional structures of actin, it clearly appears that such a method is very difficult to implement in mass production process since the distance of the two surfaces must be increased with a very precise rate, during the polymerization process, in order to promote the unidirectional growth of actin filaments. Moreover, since US 2005/0106629 does not provide any experimental result, the experimental feasibility of this method remains to be demonstrated.\nThere is thus a need for novel methods enabling the preparation of three-dimensional structures of actin on the nano- or micrometer scale while controlling precisely the shape, the orientation and the spatial location of said structures."}
{"text": "1. Field of the Invention\nThe present invention is a kind of computer architecture, a load balancing reconfigurable heterogeneous processor architecture with dynamic allocation method for high performance in particular.\n2. Description of the Related Art\nAs today's semiconductor technology advances at a rate sketched by the Moore's law, the assorted digital information apparatus tends to integrate processors with various functions into SoC (System-on-a-Chip) to suit the needs of versatility and small form factor. While such an SoC is at work, the characteristics of the application tend to use some processors of certain type intensively but leave those of other types idling from time to time, causing the abundant hardware resources often unevenly used. This ever-changing needs for different types of processors along time greatly lower the overall performance.\nFor example, the vastly used GPUs (Graphic Processing Units) in computer systems consist of large numbers of vertex shaders and pixel shaders. They process graphics through coordinate and light transformations, texture compression/decompression, bi-linear pixel shading, etc., to render graphics. The first task among these, vertex shading, shades vertices of geometries through coordinate and light transformation using a large number of vertex shaders. These shaded vertices are then passed on to another group of large number of pixel shaders and texture units for texture compression/decompression, bi-linear pixel shading, etc. As a result, often the number of pixels to be processed occasionally becomes much greater than the number of vertices, or while the vertex shaders are busy processing, the pixel shaders and texture units are idling; whereas while the pixel shaders and texture units are busy processing, the vertex shaders have little work to do. This fact makes the two sets of processors run unevenly along time, lowering the overall performance of the GPU. One solution may be to use unified shaders, but the costs are more complex shader circuits and routings.\nTo deal with such a deficiency, the US Patent US2007/0091089A1 proposes a dynamically allocateable GPU system with method, which is equipped with multiple sharable units such as a sharable vertex processor, a sharable geometry processor, and a sharable pixel processor. Through at least one control unit, the sharable processors are assigned execution tasks, and the workload of each processor is monitored. Those unloaded sharable processors can be assigned to assist the loaded sharable processors.\nHowever, the aforementioned patent US2007/0091089A1 uses a plurality of shareable shaders to share the loads of various shading tasks, resulting in complicated hardware design and its associated monitoring and load sharing algorithm. The present invention is intended to resolve such difficulties. The present invention presents dynamic reconfigurable heterogeneous processors architecture with load balancing and dynamic allocation method"}
{"text": "As an example of this type of the fluid jetting device, a local cleansing device for cleansing parts (anus for example) of the human body is well known. With this kind of the local cleansing device, upon jetting the cleansing water from a single fluid jetting spout toward a part of the human body, ordinarily, it is desirable that the water contact area of the jetted cleansing water is in some measure made to be a broad area.\nIn order to fulfill this sort of demand, a method of rotating a nozzle arm having a nozzle built therein in a circular path (nozzle arm rotation method) or a method of driving the nozzle itself within a nozzle arm having such nozzle built therein (nozzle rotation method) may be adopted. Incidentally, with the former method, since it is necessary to simultaneously control the nozzle arm in two axes on orthogonal coordinates, a drive motor or the like is required for the respective axes, and this resulted in the enlargement of the device. Meanwhile, with the latter method, since the drive target is only the nozzle, this is preferable in that the device can be miniaturized for such portion. This kind of nozzle rotation method has been variously proposed, for example, in JP 2000-8453 A, and may be classified broadly into a type that drives the nozzle with electrical power and a type that drives the nozzle with cleansing water pressure. The latter is superior to the former in terms of energy conservation.\nFIG. 1 is an explanatory diagram for explaining the structure conventionally adopted so as to rotatably drive the nozzle with cleansing water pressure, wherein FIG. 1(a) is an explanatory diagram illustrating a schematic cross section of the nozzle arm, and FIG. 1(b) is an explanatory diagram illustrating the schematic cross section of line A—A thereof.\nAs shown in FIG. 1, a truncated cone-shaped nozzle is rotatably built in the chamber of the nozzle arm, and a plurality of curved grooves is formed around the peripheral wall of the nozzle. This nozzle, at the tip side thereof, is sealed to the inner face of the chamber with a seal member. When cleansing water is supplied to this kind of nozzle, the nozzle rotates with the pressure of the cleansing water upon such cleansing water passing through the grooves between the inner face of the chamber and the peripheral wall of the nozzle. Thus, the nozzle jets cleansing water from the jetting spout at the nozzle tip so as to broaden the water contact area.\nNevertheless, with the foregoing conventional structure, since a seal member lies between the nozzle tip and inner face of the chamber, the nozzle is subject to a relatively large rotational resistance from the seal member during the rotation thereof.\nRotational speed of the nozzle affects the broadening of the cleansing water from the jetting spout, and a certain degree of rotational speed is required in order to broaden the water contact area. As a result, the water pressure upon supplying cleansing water must be increased in order to elicit and maintain the rotation of the nozzle, and problems such as the enlargement of the actuator of pumps or the like and increased operating costs would arise.\nThese problems are not typical to a cleansing water jetting device as represented with a local cleansing device, and, even with a fluid jetting device employed for other purposes, similar problems occur as a result of the structure of rotating the nozzle with fluid pressure.\nThe present invention was devised in view of the foregoing problems, and an object thereof is to seek, upon adopting the structure of rotating the nozzle with fluid pressure, the miniaturization of the actuator of a pump or the like for supplying fluid to the chamber, and the reduction of operating costs."}
{"text": "The present disclosure relates generally to medical fluid systems and more particularly to dialysis systems that can be used in a patient's home.\nIt is expected that the power requirements during a typical usage cycle of a home dialysis device will vary significantly. During therapy, which can last for two to eight hours, energy usage is expected to be particularly high. The dialysis delivery and water purification systems require significant amounts of energy to run heaters, pumps, valves, etc. Similarly, disinfection equipment and processes that are operated at the end of therapy can also require large amounts of electricity. The instantaneous amount of energy required during those states may exceed the wattage which can be supplied by a standard home electrical branch (e.g., 15 ampere, 110/120 volt, or equivalent). The patient's branch power may need to be upgraded to a higher operating current and/or voltage to run a therapy.\nIt is desirable not to have to upgrade the patient's electrical system."}
{"text": "The use of implantable medical devices for electrical stimulation of electrically excitable tissue is well known in the medical arts. For example, electrical stimulation of the brain, the spinal cord, or a peripheral nerve may be used to treat any number of medical conditions. In such devices, electrodes deliver the stimulation of signal to the electrically excitable tissue. The electrodes are operatively connected to an implantable pulse generator (IPG) that is packaged in a case that is adapted to be implantable. Those electrodes are coupled to that pulse generator by a conductive lead wire.\nA user having such an implantable medical device during normal life activities may be forced to go through a time-alternating electromagnetic field. Prevalent examples of sources of electromagnetic field are Electronic Article Surveillance (EAS) systems and Metal Detectors (MD). Such systems detect theft of articles that have an attached electromagnetic tag and are found in the exit doorways of many stores and libraries.\nEAS/MD systems work by generating an electromagnetic field that retail customers must pass through at the entrance and/or exit from a business establishment or other protected area. EAS systems detect the presence of antitheft tags on merchandise as they pass through the electromagnetic field if they have not been deactivated. Similarly, metal detectors (MD) generate an electromagnetic field that is perturbed by the presence of metal objects that might be carried by a person passing through the metal detector. Both these systems can induce voltages, via Faradays Law, on the lead system of implantable stimulators. The induced voltage is proportional to the area formed from one end of the lead to the other end of the lead (e.g. from the IPG case to the end of the lead). More specifically, according to Faraday's law, the voltage induced on a wire (lead system) is proportional to: (A) the area X number of turns described by the wire; (B) the flux density of the electromagnetic field; C) the frequency of the magnetic field; and (D) the cosine of the angle between the electromagnetic field and the vector normal to the area described by the wire. For dual lead systems, an induced voltage may also appear across the distal ends of the leads and is proportional to the area formed between the two leads.\nThese induced voltages may be of sufficient voltage and pulse duration to cause undesired tissue stimulation in the patient. This may result in shocking sensations to the patient, pacing of the heart, IPG sensing irregularities or other undesired stimulation effects. Additionally, in some medical devices with sensing, the induced voltage may also cause a false “sensing” response.\nOne technique for minimizing the effects of induced voltages is to keep the lead length of the implanted system to a minimum and/or to keep the lead in a straight line to minimize the loop area generated by the lead system. However, it may not always be possible to minimize the loop area or lead length. Another technique for minimizing the effects of induced voltages is to restrict the IPG to bipolar use only since this disconnects the case from the lead. However, this is not entirely effective since the electronic switch to the case may become forward biased (conductive) if the induced voltage is sufficiently high. Another problem with this approach is that it requires feedthrough capacitors on each of the IPG outputs to reduce the susceptibility to high frequency RF. Feedthrough capacitors can act as a low impedance connection between the lead and the IPG case at EAS/MD frequencies, thereby negating the benefit of bipolar stimulation. See U.S. Pat. Nos. 5,751,539 and 5,905,627.\nIt is therefore desirable to provide a way to improve electromagnetic compatibility (EMC) with EAS/MD systems that overcomes the disadvantages of the prior art. Moreover, it is desirable to improve EMC with EAS/MD systems while still allowing the use of unipolar stimulation."}
{"text": "My new variety of lily plant originated as a seedling which first flowered in Woodland, Wash., in 1987. The breeding efforts had as their objective the production of large-flowered completely upfacing Oriental hybrids in a variety of shades of pink, with deeply pigmented papillae, suited to forcing into flower out of season, heretofore unknown in the lily breeding art.\nI achieved the desired objective by intercrossing unnamed upright and semi-upright Oriental seedlings which carried pink flower color and were also suited to forcing for year-round use as cut-flowers. The unnamed Oriental seedlings which were the parents of `La Claridad` were produced by me, from a series of crosses among seedlings which were selected in my fields from large seedling populations of my own crosses. The original crosses came from material unique to my own breeding lines and not available in the trade. The field-selected seedlings were genetically very diverse, and 100 seedlings were selected for their outstanding vigor, stem length, pink flower color, and upfacing or semi-upfacing flower orientation. These were forced into flower under glass, and the 25 seedlings which showed the best forcing performance were then randomly intercrossed to produce a large quantity of seed. This seed was grown to flowering size, producing a large population of seedlings, from which the best 100 were again selected for vigor, stem length, pink flower color, and upfacing or semi-upfacing flower orientation. These were subsequently tested for their forcing performance, and `La Claridad` was selected as one of the seedlings best suited for forcing as a cut-flower.\nThe flowers of my new lily are characterized by a completely upfacing orientation, large size, unusually thick substance, distinctive light pink coloration with noticeably deeper pink tepal midribs, deep pink papillae, inconspicuous white-edged green nectaries, and lightly ruffled tepal margins, a combination unique among Oriental hybrid lilies. The variety possesses unusually strong, stout stems. In addition, the clone possesses to a high degree desirable characteristics of hybrid vigor. The clone is a good grower and propagator, as observed at Woodland, Wash.\nMy new variety of lily plant has been asexually reproduced by me and under my direction at Woodland, Wash. Successive generations produced by natural propagation from bulblets, by bulb scale propagation, and by tissue culturing from bulb scale explants have demonstrated that the novel and distinctive characteristics of my new variety are fixed and hold true under asexual propagation from generation to generation."}
{"text": "1. Field\nThe present invention relates to the field of heat management of computing devices, and in particular the cooling of heat generating components and exterior walls of mobile computing devices.\n2. Discussion of Related Art\nHeat management can be critical in many applications. Excessive heat can cause damage to or degrade the performance of mechanical, chemical, electric, and other types of devices. Heat management becomes more critical as technology advances and newer devices continue to become smaller and more complex, and as a result run at higher power levels and/or power densities.\nModern electronic circuits, because of their high density and small size, often generate a substantial amount of heat. Complex integrated circuits (ICs), especially microprocessors, generate so much heat that they are often unable to operate without some sort of cooling system. Further, even if an IC is able to operate, excess heat can degrade an IC's performance and can adversely affect its reliability over time. Inadequate cooling can cause problems in central processing units (CPUs) used in personal computers (PCs), which can result in system crashes, lockups, surprise reboots, and other errors. The risk of such problems can become especially acute in the tight confines found inside mobile computers and other portable computing and electronic devices.\nAs the processing powers of mobile computing devices continue to increase, the temperatures of the outer walls of the mobile computing devices will continue to rise to unacceptable levels. The temperatures are becoming the highest within the regions of the memory, central processing unit (CPU), chipset and voltage regulator (VR). To overcome the increase of heat in these locations, vents have been placed in strategic locations to reduce the temperatures.\nPrior methods for dealing with such cooling problems have included using simple vent systems in the outer walls of a mobile device. But, as the amount of cooling air available within mobile computing devices is reduced as the mobile devices are scaled down, the vent system becomes less and less efficient."}
{"text": "Brushless direct current (BLDC) motors are becoming more prevalent in industries that typically did not use BLDC motors. For example, the need for increased efficiency in the heating, ventilating, and air conditioning market has led to the use of BLDC motors for powering the blower. BLDC motors can include a rotor having a plurality of magnetic poles (e.g., a plurality of poles produced with permanent magnets) of alternating polarity disposed on a surface of a rotor core, and a stator that receives electrical power and produces a magnetic field in response thereto. The magnetic field of the stator interacts with a magnetic field of the rotor to cause movement of the rotor.\nStandard operating conditions for heating, ventilating, and/or air conditioning systems (collectively referred to herein as HVAC systems) generally vary over relatively short periods of time. It is typically preferred to adjust the cooling and heating cycles, among other parameters, of the HVAC systems as conditions vary. For example, some thermostats are configured to generate a signal indicative of a cooling requirement. In response to the signal generated by the thermostat, the motor of the HVAC system operates a fan or a blower.\nConventional HVAC applications often utilize multi-tapped permanent split capacitor (PSC) type motors. In general, a multi-tapped PSC motor is a motor that has a multi-tapped main winding where all or part of the main winding is coupled in parallel with an auxiliary starting winding that is coupled in series with a capacitor. Such multi-tapped PSC motors are used in HVAC applications, such as furnace blower and air handler applications, because the multi-tapped winding can produce variable output torque and, therefore, variable output speed for the purpose of delivering different amounts of air flow for different applications. For example, one tap setting may provide a relatively low amount of air flow for air circulation when there is no heating or cooling activity. Another tap setting could increase the air flow when cooling is desired. By using multiple taps, various operating states can be established for a tapped PSC motor, such as heating, cooling, and air. In general, each tap point on the multi-tapped PSC motor is coupled to an input line and relays are energized in response to control signals from, for example, a thermostat to energize one of the tap points at any given time.\nBlower motors used in residential HVAC systems traditionally use PSC motors. These motors generally have two independent power connections to accommodate heating or cooling modes of operation. The heating or cooling power inputs are normally connected to different winding taps in the PSC motor to provide somewhat different operating speeds for the blower in the respective modes of operation. More than two sets of taps can be designed into the PSC motor, allowing the OEM or installer to select the operating speed by appropriate connection of the taps to the respective heating and cooling power connections. The energizing of these AC power connections to the motor is controlled by activation of a temperature switch and a relay driven from the thermostat."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor package, and more particularly, to a wafer level package having connection terminals whose side surfaces are exposed and a method of fabricating the same. The present invention also relates to a semiconductor module including a wafer level package having connection terminals whose side surfaces are exposed, and a method of fabricating the same.\n2. Description of the Related Art\nA semiconductor package serves to electrically connect inputs and outputs of a semiconductor chip to an external circuit and to protect the semiconductor chip. As electronic devices get smaller and lighter and have higher performances, small, light, economical, and highly reliable semiconductor packages are required. To satisfy these requirements, wafer level packages have been developed in which the assembling or packaging of a semiconductor chip is completed at a wafer level. In the wafer level package, all semiconductor chips on a wafer are collectively subject to fabrication processes such as assembling processes, resulting in greatly reduced semiconductor device fabrication costs, better incorporation of packaging functions and semiconductor chip functions, improvements in thermal and electrical characteristics of the semiconductor devices, and miniaturization of the package to the relative dimension of the semiconductor chip.\nFIG. 1 is a cross-sectional view illustrating a conventional wafer level package 100.\nReferring to FIG. 1, a metal pad 120 is formed on a front surface of a wafer 110. A metal wiring layer 150 is electrically connected to the metal pad 120 and to a solder ball 170. A first insulating layer 130 and a second insulating layer 140 are formed between the metal pad 120 and the metal wiring layer 150, and a third insulating layer 160 is formed on the metal wiring layer 150. The first to third insulating layers 130, 140, and 160 include openings 135, 145, and 165, respectively.\nIn the conventional wafer level package 100, the metal wiring layer 150 is made of, for example, copper, and generally needs to be redistributed on the semiconductor chip in order to adhere the solder ball 170 to the metal pad 120 of the semiconductor chip for a proper electrical connection with an external circuit. Such a redistribution process of the metal wiring layer 150 increases the fabricating cost and the metal wiring layer 150 may be separated from the surface of the semiconductor chip, thereby degrading the reliability of the package.\nFurthermore, the area of the semiconductor chip limits the number of the solder balls 170 that can be arranged on the semiconductor chip. Accordingly, it is difficult to lay out a large number of solder balls 170 on the semiconductor chip. To solve this problem, the solder balls 170 can be reduced in size so that a plurality of the solder balls 170 can be more easily arranged on the semiconductor chip. However, the reduced size of the solder ball weakens the adhesion of the solder ball to the underlying metal pad."}
{"text": "In today's society, individuals and businesses conduct an ever-increasing amount of activities on and over computer systems. These computer systems, including proprietary and non-proprietary computer networks, are often storing, archiving, and transmitting all types of sensitive information. Thus, an ever-increasing need exists for ensuring data stored and transmitted over these systems cannot be read or otherwise compromised.\nOne common solution for securing computer systems is to provide login and password functionality. However, password management has proven to be quite costly with a large percentage of help desk calls relating to password issues. Moreover, passwords provide little security in that they are generally stored in a file susceptible to inappropriate access, through, for example, brute-force attacks.\nAnother solution for securing computer systems is to provide cryptographic infrastructures. Cryptography, in general, refers to protecting data by transforming, or encrypting, it into an unreadable format. Only those who possess the key(s) to the encryption can decrypt the data into a useable format. Cryptography is used to identify users, e.g., authentication, to allow access privileges, e.g., authorization, to create digital certificates and signatures, and the like. One popular cryptography system is a public key system that uses two keys, a public key known to everyone and a private key known only to the individual or business owner thereof. Generally, the data encrypted with one key is decrypted with the other and neither key is recreatable from the other.\nUnfortunately, even the foregoing typical public-key cryptographic systems are still highly reliant on the user for security. For example, cryptographic systems issue the private key to the user, for example, through the user's browser. Unsophisticated users then generally store the private key on a hard drive accessible to others through an open computer system, such as, for example, the Internet. On the other hand, users may choose poor names for files containing their private key, such as, for example, “key.” The result of the foregoing and other acts is to allow the key or keys to be susceptible to compromise.\nIn addition to the foregoing compromises, a user may save his or her private key on a computer system configured with an archiving or backup system, potentially resulting in copies of the private key traveling through multiple computer storage devices or other systems. This security breach is often referred to as “key migration.” Similar to key migration, many applications provide access to a user's private key through, at most, simple login and password access. As mentioned in the foregoing, login and password access often does not provide adequate security.\nOne solution for increasing the security of the foregoing cryptographic systems is to include biometrics as part of the authentication or authorization. Biometrics generally include measurable physical characteristics, such as, for example, finger prints or speech that can be checked by an automated system, such as, for example, pattern matching or recognition of finger print patterns or speech patterns. In such systems, a user's biometric and/or keys may be stored on mobile computing devices, such as, for example, a smartcard, laptop, personal digital assistant, or mobile phone, thereby allowing the biometric or keys to be usable in a mobile environment.\nThe foregoing mobile biometric cryptographic system still suffers from a variety of drawbacks. For example, the mobile user may lose or break the smartcard or portable computing device, thereby having his or her access to potentially important data entirely cut-off. Alternatively, a malicious person may steal the mobile user's smartcard or portable computing device and use it to effectively steal the mobile user's digital credentials. On the other hand, the portable-computing device may be connected to an open system, such as the Internet, and, like passwords, the file where the biometric is stored may be susceptible to compromise through user inattentiveness to security or malicious intruders."}
{"text": "1. Field of the Invention\nThe present invention relates to a dot-matrix type display device capable of performing large-scale display.\n2. Description of the Related Art\nA display device using field-effect liquid crystal is effective as a large-scale display device. Particularly, a display device of super-twisted nematic (STN) liquid crystal in which a structure having a large twist angle is given to the liquid crystal is known as a display device using a time-division driving characteristic. The principle of time-division driving in the above dot-matrix type display device will be described hereunder. As shown in FIG. 1, the dot-matrix type display device is provided with linear Y electrodes (signal electrodes) formed on a lower electrode substrate (not shown) and linear X electrodes (scanning electrodes) formed on an upper electrode substrate (not shown), so that display of characters, graphics, etc., is performed through an electric field applied to a field-effect liquid crystal sandwiched between the X and Y electrodes. That is, the display is performed by selective turning-on (or turning-off) of the liquid crystal at points of intersection between the X and Y electrodes.\nScanning is sequentially and repeatedly performed on the scanning electrodes X.sub.1, X.sub.2, . . . X.sub.n in FIG. 1 sequentially line by line to thereby perform timedivision driving. When, for example, the scanning electrode X.sub.3 is selected, display signals from the signal electrodes Y.sub.1, Y.sub.2, . . . Y.sub.n are applied to all pixels Z.sub.1, Z.sub.2, . . . Z.sub.n on the electrode simultaneously so as to perform selection (or non-selection). Points of intersection can be thus turned on (or turned off) by combinations of voltage pulses applied to the scanning electrodes X and to the signal electrodes Y. In this case, the number of scanning electrodes X corresponds to the number of time divisions (the number of scanning lines or the number of display lines).\nThe number of display lines is, however, naturally limited due to the characteristic of the liquid crystal. For the purpose of increasing the number of display lines, stripe electrodes (X, Y) may be separated into two groups to form two systems, which makes it possible to increase the number of display lines by twice, as typically shown in FIG. 2. For example, 800 display lines can be provided even in the case of a liquid crystal device having the number of scanning lines limited to 400. It is, however, apparent from FIG. 2 that the signal electrodes Y are connected to external driving circuits 2, 2'and the scanning electrodes X are connected to the external scanning circuits 3, 3'through end portions of the substrate 1. Accordingly, it has been impossible to separate the signal electrodes into three or more groups.\nAs a measure to solve the aforementioned problem there has been proposed a method of separating the scanning electrodes into three or more groups in which holes are formed so as to substantially perpendicularly pass through the substrate on which the signal electrodes are formed, and in which electrical conductive leads are inserted into the through holes respectively, connected to the signal electrodes and led out to the back surface side of the substrate.\nAccording to the proposed method, the signal electrodes can be separated into a desired number of groups, so that the total number of display lines can be increased. In this method, however, occurrence of irregular display called \"crosstalk\" cannot be suppressed perfectly. The above \"crosstalk\" means a phenomenon that the display contrast ratio of the whole display screen is reduced because of occurrence of imperfect turning-on at non-display points (non-selective pixels) not to be displayed. This is a weak point of the dot-matrix display device. When, for example, a certain driving voltage is set, the liquid crystal at non-selection points responds, even through imperfectly, to the driving voltage to increase transmittivity because a bias voltage is applied also to the non-selection points. That is, the difference between the transmittivity in this state and the transmittivity in the case of attaining a perfect non-selection state is called \"crosstalk\". The crosstalk phenomenon changes depending on the number of time divisions, the electrode resistance and the connection resistance. Generally, crosstalk is apt to occur as the number of time divisions increases, and it becomes large as the sharpness of the threshold characteristics of the voltage and transmittivity becomes low as the electrode resistance and the connection resistance become large."}
{"text": "As disclosed in a patent document, Japanese Patent Laid-Open No. JP 2004-326629 A, (patent document 1) listed below, an abnormality monitoring device that monitors an abnormality of the microcomputer is known. The abnormality monitoring device includes an abnormality detector that detects an abnormality during a normal operation of the microcomputer and a standby disorder detector that detects a disorder during a standby time of the microcomputer.\nThe microcomputer outputs a watch dog pulse during the normal operation, and does not output the watch dog pulse during the standby time. Therefore, the abnormality detector determines that the microcomputer in the normal operation has an abnormality caused therein when the watch dog pulse is not input to the abnormality detector from the microcomputer for more than a preset period of time. The standby disorder detector determines that the microcomputer has a disorder caused therein when the watch dog pulse is input to the standby disorder detector from the microcomputer during the standby time of the microcomputer.\nAs described above, the standby disorder detector disclosed in the patent document 1 determines that the microcomputer has a disorder caused therein when the watch dog pulse is input to the standby disorder detector from the microcomputer during the standby time of the microcomputer. However, depending on a state of disorder of the microcomputer in the standby time, the microcomputer may possibly not output the watch dog pulse during the standby time (i.e., designated as a low power consumption mode, or a low power mode). Therefore, the abnormality monitoring device in the patent document 1 is not capable of detecting such a disorder of the microcomputer."}
{"text": "Generally, most organizations that utilize computer systems and networks monitor those systems and networks for changes. This monitoring is often initiated and performed by a server that requests a number of client settings from monitored computing devices, servers, peripheral devices, etc. within the system or network. Upon receiving the settings, the server may then analyze, classify, and/or store them, and issue a report that indicates changes to the settings. Often, such a report includes numerous occurrences of various changes. Indeed, the number of occurrences may be so numerous that it renders the report nearly useless."}
{"text": "Recently, the global warming phenomenon by carbon dioxide has become an issue, and to reduce the carbon dioxide emission amount has become an assignment. In automobile as well, to reduce the carbon dioxide amount in the exhaust gas has become an assignment, and a lean-burn engine which can thinly burn a fuel in an excessive oxygen atmosphere has been developed. By this lean-burn engine, since mileage is improved, it is possible to suppress the emission amount of carbon dioxide.\nBy the way, in a case where harmful components in an exhaust gas from a lean-burn engine are reduced, since the exhaust gas is put in an excessive oxygen atmosphere, reduction reactions become less likely to occur, and the reduction and purification of NOx become difficult. Hence, in Japanese Unexamined Patent Publication (KOKAI) No. 5-317,652, an NOx storage-and-reduction type catalyst is disclosed on which an NOx storage component, being selected from the group consisting of alkali metals, alkaline-earth metals and rare-earth elements, is loaded along with a noble metal. By using this NOx storage-and-reduction type catalyst and controlling an air-fuel ratio so that it becomes from a fuel-stoichiometric to rich atmosphere in a pulsating manner in the midway of a fuel-lean atmosphere, it is possible to efficiently progress the oxidation of HC and CO as well as the reduction of NOx, and accordingly it is possible to acquire high purifying performance.\nNamely, an exhaust gas which is burned in a fuel-lean atmosphere becomes a reduction components lean atmosphere, in the reduction components lean atmosphere, NO in the exhaust gas is oxidized to become NOx and is stored into the NOx storage component, and accordingly the emission of NOx is suppressed. Then, when it is controlled from a fuel-stoichiometric to rich air-fuel ratio in a pulsating manner, the exhaust gas becomes from a reduction components stoichiometric to rich atmosphere. Therefore, NOx are released from the NOx storage component, and they react with the reduction components, such as HC, which exist in the exhaust gas, so that they are reduced, and accordingly the emission of NOx is suppressed. Therefore, it is possible to suppress the emission of NOx in all of the atmospheres from fuel-rich to fuel-lean.\nTo control an air-fuel ratio in a pulsating manner so as to become from a fuel-stoichiometric to rich atmosphere is referred to as rich spiking, and the extent of making a fuel-rich atmosphere by the rich spiking is expressed by deep or shallow. Namely, by rich spiking, making a heavy degree of fuel-rich atmosphere is referred to as “charging rich spiking deeply,” and making a light degree of fuel-stoichiometric to rich atmosphere is referred to as “charging rich spiking shallowly.” Then, in Japanese Unexamined Patent Publication (KOKAI) No. 11-107,810, and the like, there is set forth to appropriately control the extent of the rich spiking and the timing of charging it.\nThen, it has been required to purify NOx, which are emitted from an engine being operated under a variety of fuel-lean conditions, by using an NOx storage-and-reduction type catalyst while always sustaining a purifying rate as high as 90% or more. However, in conventional NOx storage-and-reduction type catalysts, when the NOx amount, which is emitted from an engine in a unit period of time, is large, the rich spiking should be charged at intervals of from a couple of seconds to 60 seconds. However, in the excessive fuel component to be supplied as the rich spiking, there exists, in addition to an amount to be used in the reduction of NOx, an amount to be used in controlling the combustion state of an engine to a fuel-rich air-fuel ratio. Therefore, in a case where the rich spiking is charged frequently, there arises the deterioration of mileage. In particular, in a case where being driven ordinarily at a high speed, the NOx amount, which is emitted from an engine, enlarges remarkably. Therefore, in order to sustain a high NOx purifying rate, the rich spiking should be kept being charged at very short intervals of 10 seconds or less, and accordingly there is a problem in that the mileage is lowered considerably.\nThe present invention has been done in view of such circumstances, and its main object is to improve mileage by prolonging the intervals of the rich spiking as well as to improve the reduction-and-purification efficiency of NOx.\nThe characteristics of an exhaust gas purifying apparatus for an internal combustion engine, which solves the aforementioned assignments, lie in that it is used for an internal combustion engine, which can select an operation at a fuel-lean air-fuel ratio and an operation at a fuel-stoichiometric or rich air-fuel ratio and comprises: an NOx storage-and-reduction type catalyst disposed in an exhaust gas flow passage and exhibiting a saturated NOx storage amount of 5 g or more as NO2 with respect to 1 liter of a catalyst volume at 500° C.; NOx storage amount estimating means for estimating an actual NOx storage amount of the NOx storage-and-reduction type catalyst; air-fuel ratio adjusting means for adjusting an exhaust gas atmosphere to reduction components lean or reduction components rich; and a controlling device for controlling the air-fuel ratio adjusting means based on an estimated value estimated by the NOx storage amount estimating means so that the actual NOx storage amount becomes 50% or less of the saturated NOx storage amount.\nMoreover, the characteristics of an exhaust gas purifying process for an internal combustion engine of the present invention lie in that, in an exhaust gas purifying process for an internal combustion engine, in which an NOx storage-and-reduction type catalyst including an NOx storage component is contacted with an exhaust gas from an internal combustion engine, which can select an operation at a fuel-lean air-fuel ratio and an operation at a fuel-stoichiometric or rich air-fuel ratio, thereby storing NOx contained in the exhaust gas into the NOx storage component in a reduction components lean atmosphere, and reducing NOx released from the NOx storage component by making a reduction components stoichiometric to rich atmosphere by rich spiking, an NOx storage-and-reduction type catalyst exhibiting a saturated NOx storage amount of 5 g or more as NO2 with respect to 1 liter of a catalyst volume at 500° C. is used, and in that the rich spiking is controlled so that an actual NOx storage amount of the NOx storage-and-reduction type catalyst becomes 50% or less of the saturated NOx storage amount.\nThe “fuel-lean air-fuel ratio” refers to an air-fuel ratio which makes an exhaust gas into an atmosphere in which oxygen exists in a concentration which exceeds an oxygen equivalent ratio required for oxidizing all of reduction components, such as CO, THC and H2. In the case of an air-fuel ratio (A/F: ratio of air to fuel by weight), the equivalent (stoichiometric) point is around 14.6, and, in the present invention, an atmosphere whose A/F exceeds 14.6 is called a fuel-lean air-fuel ratio. Moreover, the “fuel-rich air-fuel ratio” refers, contrary to the “fuel-lean air-fuel ratio,” to an air-fuel ratio which makes an exhaust gas into an atmosphere in which oxygen exists in a concentration which does not reach an oxygen equivalent ratio required for oxidizing all of the reduction components, an air-fuel ratio in which A/F does not reach 14.6 is called a fuel-rich air-fuel ratio."}
{"text": "This invention relates generally to brackets for positioning detection equipment on walls and the like and more specifically to a bracket having antitampering features."}
{"text": "With some vehicles the driver and passenger enclosure is provided by separate units which are joined together. Thus, cab and sleeper units of a truck provide one example of such a vehicle. The cab and sleeper are typically independently bolted or secured to the vehicle chassis. During operation, the truck cab and sleeper box move relative to one another as the chassis twists in response to travel across uneven ground. To accommodate this relative movement, the cab and sleeper are generally separated by a gap of up to three inches. The cab and sleeper units typically have opposed walls which are formed with an opening for convenient passenger access between the compartments. A weather seal is provided between the cab and sleeper to prevent moisture from entering the separation gap and passing into the passenger compartment when the vehicle is exposed to either elements.\nOne type of weather seal used with vehicles is a rubber boot. Mounting frames formed for retaining the rubber boot are secured to the opposed walls of the cab and sleeper units. Several attachment screws secure the rubber boot to the opposed mounting frames, substantially enclosing the separation between the adjacent units. A compression ring is often included to evenly distribute the forces which are exerted by the attachment screws. The rubber boot protects the interior of the passenger compartment by substantially preventing water and debris from passing through the aligned openings in the opposed walls, cab and sleeper box walls.\nSeveral fastener-receiving holes are formed in the boot, with the attachment screws being inserted into the holes to mount the boot to the cab and sleeper units. When the vehicle is exposed to the elements, moisture may seep around the screws and through the holes formed in the boot. The moisture may then collect within the gap and pass into the interior of the vehicle, damaging the upholstery and potentially inducing the formation of rust. During operation of the vehicle, the boot will partially absorb the forces exerted by the relative movement of the cab and sleeper units. However, the vibration of the chassis and the forces exerted on the boot by the motion of the cab and sleeper units may slightly enlarge the size of the fastener-receiving holes, increasing the amount of moisture entering the passenger compartments.\nThe forces applied to the boot and attachment screws during operation may additionally gradually loosen the screws, with the rubber boot eventually becoming separated from the mounting frames. The released portions of the boot must then be secured to the mounting frame, or if necessary the boot must be replaced. Unless the driver has access to the tools required for replacing the seal, the driver must return the vehicle to an appropriate service facility if the boot loosens or pulls free during operation. Thus, maintenance, repair and replacement of the rubber boot assembly can be a time consuming interruption which may potentially significantly delay the truck driver.\nA flexible boot assembly which is secured to the adjacent compartments without the use of fasteners and fastener-receiving holes would substantially prevent moisture and debris from entering the interior of the vehicle. By eliminating the use of screws or other types of fasteners, efficient replacement of the boot assembly may be completed without requiring specialized tools and a plurality of replacement parts. Such a fastener-free boot assembly would substantially prevent the entrance of moisture into the vehicle interior, even after long periods of continuous operation. Such a boot assembly would also absorb the forces exerted by motion of the cab and sleeper units without developing potential leakage paths.\nOne type of prior art flexible boot currently commercially available is mounted to a frame provided on the opposed walls of the vehicle compartments without attachment screws. The frame typically includes an opening formed to receive the flexible boot, which has opposed, arrowhead-shaped edges. The arrowhead-shaped edge, which is formed for insertion into the opening on the flange, is forced into the opening to secure the boot to the adjacent cab and sleeper compartments.\nThe arrowhead-shaped boot edges secure the sealing boot between the opposed walls of the adjacent compartments to prevent the leakage of water into the vehicle interior. The barb portion, large enough in size to hold the boot in place, requires a substantial installation force. The edges must be hammered or otherwise forcefully driven into the openings in the frame member. The installation of the sealing boot is both manually taxing and time consuming. The installation process is impeded by the restricted work area provided between the cab and sleeper. Moreover, to prevent release of the boot from the frame during vehicle operation, the arrowhead barb must be completely retained within the frame.\nDuring operation of the vehicle, debris, snow and ice will collect on top of the sealing boot. The operational forces and vibration of the vehicle, and the weight of the material collecting on the boot, gradually urge the arrowhead-shaped barbs from the openings of one or both of the mounting frames. The arrowhead-shaped edge must be reinserted to prevent water from leaking into the interior of the vehicle. If the boot is not replaced, moisture will enter the compartments, eventually damaging the vehicle interior. Since the sealing boot is difficult to install, the time required to replace the sealing boot will significantly detain the driver. When replacement of the boot is necessary, the arrowhead-shaped barb must be disengaged from the channel. Manually removing the arrowhead-shaped barb from the channel is at least as difficult and time consuming as the insertion process, as the barb must be forcefully removed from the channel.\nAccordingly, it would be highly desirable to have a fastener-free flexible boot assembly which may be installed without requiring the application of a significant amount of force to mount the boot to the vehicle. Should the boot assembly inadvertently become disengaged during operation, such a boot should be capable of being rapidly and efficiently re-installed without the expenditure of substantial time and effort. Moreover, such a fastener-free, flexible boot should be formed to withstand the forces applied during operation of the vehicle, including the weight of the snow and ice collecting on top of the boot, and should remain substantially intact during operation even in severe weather conditions. Still further, a flexible boot which may be manually removed without substantial effort would significantly improve the ease with which the cab and sleeper may be disassembled or the boot replaced.\nAccordingly, the present invention provides a flexible boot assembly for obstructing the passage of water into the vehicle interior. The present boot assembly is installed without the use of screws, other types of fasteners or specialized tools, and it is easy to install and manually remove. Moreover, the fastener-free boot assembly of the present invention can be used in a novel method for conveniently and efficiently installing the boot assembly.\nA primary object of the present invention, therefore, is to provide a boot system for lining the passageway between compartments of a vehicle.\nA further object of the present invention is to provide a boot assembly which may be easily installed.\nAnother object of the present invention is to provide a boot system which will withstand the operational forces of the vehicle, preventing disengagement of the boot assembly during operation of the vehicle.\nYet another object of the present invention is to provide a boot system which may be manually removed without difficulty.\nAn additional object of the present invention is to provide a boot system which will withstand the weight of materials collecting on an exterior surface of the boot assembly.\nA more general object of the present invention is to provide a boot system which may be inexpensively manufactured and efficiently and conveniently assembled.\nThe boot assembly and method of the present invention have other objects and features of advantage which will become apparent from and are set forth in more detail in the description of the Best Mode for Carrying Out the Invention and the accompanying drawings."}
{"text": "Scientists, technicians and other users often perform spectral microanalysis of specimens to determine the spatial composition of the specimen—in other words, what makes up the specimen (elements, compounds, alloys, phases, etc.) at various points about an area of interest on the specimen. In essence, spectral microanalysis yields a map of the specimen's composition, and can do so on a microscopic level, yielding an extremely useful tool for science and industry.\nSpectral microanalysis is usually performed by directing an excitation beam—such as an electron beam or X-ray beam—at a specimen, and then capturing and analyzing the radiation and/or particles emitted by the specimen in response to the excitation beam. Often, data related to the emitted radiation and/or particles are presented to a user as a spectral histogram (or simply a “spectrum”), wherein the counts (numbers) of emissions from the specimen are plotted versus their energy. Most distinct substances generate a unique spectrum, and thus the substance's spectrum serves as a “fingerprint” for the substance. By capturing spectra from a specimen and then cross-referencing these spectra versus catalogues/databases of “reference spectra”—spectra captured from known substances—one can determine the substances present in the specimen. Usefully, the aforementioned spectra can often be collected along with an image of the specimen, thereby allowing an analyst a useful visual reference against which the spectral (substance identification) data may be reviewed. Thus, one can both view the analyzed region of interest on the specimen, and also determine the composition of the specimen at the viewed region.\nSeveral factors tend to complicate microanalysis. Initially, since specimens rarely consist of pure substances, one usually cannot simply identify the “closest” reference spectrum and assume that it identifies the substance. Rather, a spectrum captured from a specimen often reflects the presence of several substances, each being present in different amounts within the specimen. Thus, the specimen spectrum equates to a combination of the spectra of the component substances, with each component spectrum having a strength within the combination which corresponds to the relative amount of that component. As a result, analyzing a spectrum captured from a specimen to determine its component spectra (and this its component substances) often requires complicated and time-consuming statistical methods.\nFurther, the method described above only provides information about (i.e., a spectrum for) the area of the specimen which is subjected to the excitation beam. Thus, if one wishes to obtain information about a larger region of the specimen, the excitation beam must be scanned across the specimen (e.g., in an X-Y “rastering” pattern), and spectra must be collected at multiple points during the scan. This leads to large amounts of data, and consequently large data analysis times.\nTo better illustrate these difficulties, it is useful to consider common spectral microanalysis methods in greater detail. As an excitation beam is scanned between areas or “pixels” distributed about a specimen, a spectrum is collected at each pixel. Each spectrum contains information regarding the elements present in the specimen at the pixel from which the spectrum was captured. In addition, data resulting from the excitation beam is also captured at each pixel and can be used to generate an image of the pixel. For example, where the excitation beam is an electron beam, the image may depict the pixel's backscattered electrons (electrons from the excitation beam which were “reflected” from the specimen), or the pixel's secondary electrons (electrons knocked out of the specimen by the excitation beam). In either case, the image provides a visual representation of the pixel, though the visual representation may not correspond to the pixel's appearance if viewed by the eye under standard light; for example, a backscattered electron image effectively provides a view of the pixel's density, and a secondary electron image effectively provides a view of the pixel's surface roughness. The combination of the spectrum and image data at a particular pixel is often referred to as a “hyperspectral image,” or as a “hyperspectral data cube” when referring to the spectra and image data over the entire scanned region (i.e., over all pixels).\nThe spectra from the various pixels are then statistically analyzed via Principal Component Analysis (PCA) to determine which components (i.e., compounds or other combinations of elements) seem most likely to be present in the scanned region of the specimen, and at each pixel in the scanned region. In PCA methods, the correlations between the pixel spectra are determined to extract the spectra of the underlying components that seem to be present in the various pixels across the scanned region. This results in a set of statistically-derived spectra, each representing a component, wherein these derived spectra of all of these (derived) components combine in varying proportions at each of the pixels to result in the measured spectra at the pixels.\nThe information regarding the derived components and their proportions can then be extended to the measured hyperspectral images to determine their likely components. For example, one could look to each pixel and determine the dominant component at that pixel (i.e., the derived component which seems present in greatest amounts at that pixel, and/or which seems to be most likely to be present at that pixel), and then combine/average the measured spectra of all pixels having that same dominant component. The combined spectrum—which at least theoretically represents the component(s) present at its respective pixels—can then be assigned to each of the pixels having that same dominant component. This step, which replaces the measured spectrum for each pixel with its corresponding combined spectrum, can result in a smaller and more manageable set of spectra for the pixels. These spectra can then each be processed versus databases of reference spectra to determine the probable real-world components at each pixel. If desired, the image(s) of the specimen can then have the name(s) of the calculated components overlaid over the areas to which they correspond, or the image can be color-coded or otherwise encoded/labeled to allow a user to visualize the composition of the specimen.\nThe aforementioned statistical analysis of the measured pixel spectra is often performed by a commercially available software package known as COMPASS, details of which are described in greater detail in U.S. Pat. Nos. 6,584,413 and 6,675,106. While COMPASS is a valuable microanalysis tool, it suffers from the drawback that it can be time-consuming to use: since there are many pixels (e.g., over a million pixels across a 1024×1024 pixel area), and the correlations between their unique measured spectra must be calculated to determine their probable components (and these spectra may themselves consist of thousands of data points, e.g., counts measured across 4096 energy intervals or “bins”), it can require many minutes to provide the desired output, even where exceptionally fast computer processing speeds are used. Other analysis packages, which may employ different statistical processing methods than those of the aforementioned patents, tend to incur similar processing times. This is disadvantageous because an analyst naturally wants microanalysis results as soon as possible."}
{"text": "1. Field of the Invention\nThe present invention relates to card sorting devices, and more particularly, to automatic card sorting devices.\n2. Statement of the Problem\nCard games requiring a plurality of game cards are quite common. Card games generally require that the plurality of cards be reordered (shuffled) to provide a random distribution so that no advantage is given to any participant in the play of the game based on the likelihood of a particular card being at or near a particular location. To achieve this random distribution the cards are usually shuffled several times by a person acting as a dealer. The processes typically employed by dealers result in reordering the game cards in a deck that approaches but does not achieve randomness. Consequently certain skilled individuals are able to locate or nearly locate cards critical to the play with a degree of accuracy sufficient to gain a sizable edge in the play of the game. This process has been well documented for the game of BLACKJACK and is normally referred to as \"shuffle-tracking\".\nRecently, mechanical shufflers have been designed to perform the shuffling, and some of these are controlled by computer components to assist in the physical movements involved. However, all current methods involve a mechanical shuffle mimicking the process used by the human dealer, and like the shuffle produced by a human dealer, such methods approach but do not achieve randomness. Many of these machines have a tendency to move a number of cards at a time so that it is possible that several adjacent cards will remain together before and after the shuffle. In any event, shuffles produced by a human or a mechanical device are difficult to test for randomness in the real environment of a gaming table.\nAdditionally, it is desirable to know that the plurality of cards is intact in terms of card suit information, the number of cards of a specific rank, as well as total number of cards in the plurality. This is valuable in that the introduction of additional cards or the removal of certain cards can have a dramatic effect on the probabilities involved in the play of the game.\nAll shuffling techniques produce a certain amount of stress or friction to the card surfaces thus affecting the life of the cards. Therefore, it is desirable to achieve the shuffle with as little card movement as possible. Mechanical designs tend to require repeated shuffles (or \"riffles\") to achieve randomness.\nIt is further desirable that a card sorting device be used to return the plurality of cards to the \"as new\" sequence (referred to herein as the \"new\" sequence) for storage, sale, or reintroduction to the game at a later time. In BLACKJACK it is common to frequently introduce new decks of cards frequently. When removed from the tables, these decks are sorted back into new sequence by hand, stored and sold in gift ships, vending machines located in the casino or to various wholesalers. After sorting, these cards are usually drilled, dog-eared, or otherwise altered so as to distinguish these from new decks. It is estimated that some 10,000 BLACKJACK tables in America use over 30 million decks each year, most being sorted back to original (\"new\") order by hand.\nFurthermore, several games require a specific sequence that is not random, but is different than the sequence of a typical \"new\" deck. For instance, in poker, the deck is presented to the table for play by spreading the cards on the table. The suit order required in poker is spades, hearts, clubs, and diamonds. It is common to use cards that are alternatively presented to the table for play, removed and resorted to this sequence, and then used in play again. Such a sequence, which is not random but dictated by the play of the game, is referred to herein as a \"game play\" sequence. The constant sorting is currently accomplished by various employees as a secondary duty.\nIt is also desirable to verify that a winning hand indeed came from the dealer's deck and did not come from another deck. Many games offer large prizes for certain infrequent card combinations, such as Caribbean Stud in which the jackpots are often over $100,000. The added security of being able to verify a hand by checking the card sequence of the card deck is of importance to casinos offering such games in order verify winning hands.\nIt would be desirable to have a single device to solve the above stated problems that is compact, adaptable to mounting on a gaming table, and that requires little operator skill or training. Such a device should shield the sorting process from view of the players, and should be free of annoying sounds or vibrations. Also, the device should be fast, accomplishing a complete sort in less than one minute. The prior art is devoid of any single device or process capable of sorting into a random, new, or game play predetermined sequence and having these characteristics.\n3. Solution to the Problem\nThe present invention satisfies the aforementioned needs by providing an automatic card sorter that either provides a shuffled stack with a predetermined random sequence or sorts a stack into a new, game play, or other user-defined predetermined sequence. All (random, new, game play, or other user-defined) sequences are referred to herein as \"predetermined\" because the sequence of the final sorted deck is determined before any card is moved. Although the present invention can be used to sort a large variety of cards, the present discussion, for illustration purposes, will discuss the invention in the context of a stack of one or more standard decks of game cards, each card having a rank and suit printed on its face.\nThe automatic card sorter of the present invention includes at least one card reader for reading the characters on the face of a game card. From the input from the card reader, the controller, typically a microprocessor in conjunction with appropriate software and hardware, determines the rank and suit of the card. The controller also has sufficient memory to store a variety of information. The card reader, in an alternative embodiment, can also read identification codes (other than rank and suit information).\nHaving been given a stack to be sorted through an appropriate user interface, the controller provides a predetermined sequence for the cards. This predetermined sequence can be random as provided by a random number generator, the same sequence as a new deck (the \"new\" sequence), or game play (such as required by poker described above) predetermined sequence. An operator chooses the desired predetermined sequence, either a random sequence, \"new\" sequence, or game play sequence. Additional user-defined predetermined sequences could be provided for. Once the choice is made, the controller provides the predetermined sequence by retrieving the appropriate sequence from either a random number generator (when the choice is \"shuffle\") or memory.\nThe automatic sorter also includes a moveable tray having at least as many slots as cards in the stack. In the preferred embodiment, the first slot of the tray will receive the card having the value of the first position in the predetermined sequence. Hence, once the predetermined sequence is selected, each of the slots of the tray has an associated value (rank and suit) corresponding to the desired predetermined sequence. In operation, each card is read for its value. The controller then causes the tray to be moved to align the slot with the value of the read card into a position to accept the read card. The controller then causes the read card to be moved into the slot by a card feed mechanism consisting of a series of rollers in the preferred embodiment. The next card is then read, and the tray is moved to align the slot assigned to this card, which is then moved into the slot. This process is repeated until all of the cards in the stack have been placed in the tray in the predetermined sequence. The tray is then rapidly spun so that the cards are removed from the tray to form a sorted stack with the cards in the predetermined sequence, which can then be removed.\nThe random number generator, when the random mode (or \"shuffle\") for the predetermined sequence is selected, assures that the shuffled deck is random. Preferably, a random number generator such as that currently used in electronic poker machines and that has been approved by a gaming commission is used, so no additional gaming commission approvals will be required. The electronic random number generator and subsequent sorting into the sequence generated provides a degree of randomness that is beyond the capability of mechanical processes and is universally approved by gaming commissions. By reading the cards as they are sorted into the electronically generated random sequence, the randomness can be demonstrated by comparing the sequence of the cards in the stack after the shuffle to the generated random number as well as to the original unsorted stack.\nAs stated above, one reason that mechanical shufflers and human dealers fail to achieve randomness is that several cards are often moved together and may remain adjacent through several shuffles (\"riffles\"). The prior art shufflers cannot detect when multiple cards have been moved, so the possibility is always present that several cards could remain together through out a shuffle. The preferred embodiment of the automatic card sorter counts each card as it is read. When a smaller number of cards is counted, the possibility exists that the sorter mistakenly moved two or more cards at once into the same slot of the tray. While the tray and mechanism for moving the cards have been designed to move only one card at a time, the sorter further assures this result by counting the cards read. If fewer than the expected number of cards is counted, the automatic card sorter will alert the operator by appropriate output means, and the stack can be removed and replaced.\nBecause each card in the stack is read, the controller can store in memory the sequence of the original stack of cards as well as count the cards in the stack. By counting and reading the cards, the present invention verifies that the correct number of cards having proper rank and suit is present, thereby alleviating the possibility that cards have been removed, added, or substituted (i.e., replacing \"5's\" with face cards) during previous play. Once the operator inputs the number of decks present in the stack, the sorter can recall from memory the correct number of each rank and suit that should be present in the stack. While reading the cards in the stack, the controller can keep track of the cards read and alert the operator through appropriate output means of the low or high count. If fewer or more than a certain suit or rank are found, the controller alerts the operator through appropriate output devices that the stack has been subjected to tampering thus verifying the integrity of the deck and eliminating one of the favorite methods of card cheats.\nThe automatic card sorter of the present invention also eliminates wear to the surface of the playing cards. Only one \"pass\" is required in contrast to the many \"riffles\" used by human dealers or multiple shuffles required by mechanical shufflers imitating human motions. The automatic card sorter moves each card only once achieving true randomness, whereas the prior art methods require several movements resulting in greater wear while not achieving complete randomness. When such randomness is predetermined by a computer and achieved by reading and then relocating the card, the wear on the card surface is reduced accordingly because only one movement of each card per shuffle is required, whereas with a dealer or mechanical devices, several \"riffles\" are often conducted in an attempt to achieve randomness.\nAn additional advantage of the automatic card sorter over the prior art mechanical shufflers is that it can store in memory a series of predetermined sequences. This series could be recalled later to show, for example, the sequence of every sleeve or stack dealt at a given table where the automatic card sorter is installed. This information, when combined with other devices that keep track of the players at a gaming table and the hands dealt to those players, could be used to verify the winning hands at the table.\nThus, a single device is provided that is compact, can be mounted on a gaming table, and requires little operator skill or training. To operate, the operator simply places a stack of cards in a deck holding area, through which the cards are presented to the card reader (or read heads), inputs the number of decks in the stack, and selects the desired predetermined sequence. The tray and card feed mechanism are enclosed thereby shielding the entire sort from the view of the players and dealer. Thus, the present invention provides a single device and process capable of sorting into a random, new, game play, or other predetermined sequence that, in the preferred embodiment, is fast and free of annoying sounds or vibrations."}
{"text": "1. Field of the Invention\nThe invention relates to a signal delay device, and more particularly to a digital delaying device.\n2. Description of the Related Art\nThere are various formats in the data write control specifications for an optical disc such as CD-R, CD-RW, DVD, and the like. The formats, which allow writers to write multilevel data to the optical disc, are sequentially proposed. Therefore, it is more important to precisely control the laser waveform for writing data to the optical disc.\nFIG. 1 shows a schematic illustration of a laser write waveform of a CD-R. As shown in FIG. 1, the laser write waveform is a write power from time t0. The write power is held until time t1 and is dropped to an erase power, which is held until time t2. Next, the laser write waveform is again dropped to a read power. The length from time t1 to time t2 is determined by the mark length to be written. The length from time t0 to time t1 is program-controlled and does not vary with the change of the mark length. The basic unit of the mark length is T, and the mark length ranges form 3T to 11T in the CD-R specification. In addition, the leading edge of the laser write waveform at time t0 has about a delay length Δt, which is substantially equal to 0.25T. The delay length Δt is only needed when the previous write mark is 3T. Consequently, it is quite important to precisely control the switch timing of the laser write waveform.\nU.S. Pat. No. 6,269,060, entitled “Programmable write signal generator” has disclosed a technology for generating a write signal. FIG. 2 shows a block diagram of the delay signal generating unit of this patent. Referring to FIG. 2, the delay signal generating unit 506 includes an upper course delay unit and a lower fine delay unit. The output signal of the course delay unit is supplied to the fine delay unit, which may precisely control the delay timing. As shown in the drawing, the delay timing of the course delay unit is 2.5 ns while the delay timing of the fine delay unit is 0.25 ns. Each delay unit includes a plurality of cascaded delay cells D1(D2), a multiplexer 612(632), and a selection signal generator 610(630). The plurality of cascaded delay cells D1(D2) receives the input signal, and then a plurality of signals with different delay timings are generated. Next, the multiplexer 612(632) selects one of the signals for output according to a selection signal, which is supplied from the selection signal generator 610(630). Because the patent utilizes voltage control delay cells in the delay lock loop (DLL), it discloses an analog control method, and the charge pumps in the DLL are also of an analog design. In general, the analog design is more complicated and tends to be influenced by the environment."}
{"text": "The present invention relates to signal lamps, and relates more particularly to a neon tube type signal lamp assembly for use with the gear lever of a motor vehicle.\nFIG. 1 shows a neon lamp device installed in the gear lever of a motor vehicle. This structure of neon lamp device has an ON/OFF switch mounted on the outside of the transparent shell thereof and controlled to turn on/off the neon tube inside. The internal electric circuit of the neon lamp device is extended out a hole on the transparent shell and inserted through a hole on the bottom covering of the gear lever and then connected to the electric circuit of the motor vehicle. Therefore, the neon lamp device cannot be detached from the gear lever for a separate use when installed."}
{"text": "For a long period of time television sets have been designed with a 4:3 aspect ratio. That is, the width of a television set is 4/3 times its height. The 4:3 aspect ratio is often referred to as the standard format. In the United States, these sets are able to receive National Television System Committee (NTSC) signals and the video data encoded in NTSC signals are in the 4:3 aspect ratio. Most of the television shows produced are recorded in a 4:3 aspect ratio.\nMovies produced for theater exhibition use a widescreen format. The most common widescreen format utilizes a 16:9 ratio. High definition television sets are often produced in a 16:9 ratio instead of the 4:3 ratio. High definition television signals also carry video data that are in a 16:9 ratio. A number of television shows are now produced using a widescreen format.\nThese high definition television sets are also capable of displaying standard video. However, because there is extra horizontal space on a television with a 16:9 ratio compared to a 4:3 output, black vertical bars are often added to the left and right of the displayed video. By stretching a 4:3 image horizontally to fully utilize a 16:9 display, the result would often be people and other objects appearing unusually short and stocky."}
{"text": "The present invention relates generally to operations performed and equipment utilized in conjunction with subterranean wells and, in an embodiment described herein, more particularly provides a system of measuring properties and transmitting measurements in wells.\nIt is very beneficial to be able to measure properties, such as pressure, in a well and then transmit those measurements to a remote location, such as the earth\"\"s surface or another location in the well. For example, a valve system described in U.S. Pat. No. 6,152,232, the entire disclosure of which is incorporated herein by this reference, permits a drill string to be conveyed through a casing string in an underbalanced condition. The system includes a valve which is opened when the drill string is run into the casing string, and the valve is closed when the drill string is retrieved from the casing string.\nIn order to open the valve, it is desirable for there to be no pressure differential across a flapper of the valve, or at least for the pressure differential to be known before opening the valve. Unfortunately, however, there presently exists no satisfactory means for measuring this pressure differential at the time the drill string is run into the casing string, or for transmitting the pressure differential measurements to a remote location. Therefore, an operator must estimate the pressure differential by making certain assumptions, calculating hydrostatic pressure at the valve, etc., which leads to errors in the pressure differential estimate.\nIn carrying out the principles of the present invention, in accordance with an embodiment thereof, a system and method are provided which solve the above problems in the art, and which are useful in other situations, as well. In an example described below, at least one sensor is used to sense a pressure differential across a closure member of a valve. An indication of the pressure differential is transmitted from the valve to a tool positioned in a flow passage of the valve for transmission to a remote location.\nIn one aspect of the invention, a method of measuring a pressure differential across a closure member of a valve in a well is provided. The method includes the steps of: positioning the valve in the well, a flow passage extending longitudinally through the valve, and the closure member blocking flow through the flow passage; sensing the pressure differential in the flow passage across the closure member using at least one sensor positioned in the valve; conveying a tool into the flow passage; and transmitting an indication of the sensed pressure differential from the valve to the tool.\nIn another aspect of the invention, a valve for use in a well is provided. The valve includes: a flow passage formed through the valve; a closure member operative to selectively permit and prevent flow through the flow passage; and at least one sensor operative to sense a pressure differential in the flow passage across the closure member.\nIn yet another aspect of the invention, a valve system for use in a well is provided. The valve system includes: a valve positioned in the well, the valve including a closure member selectively permitting and preventing flow through a flow passage formed through the valve, and at least one sensor operative to sense a pressure differential in the flow passage across the closure member; and a tool positioned in the flow passage, the tool transmitting power to the valve to operate the sensor, and the valve transmitting an indication of the sensed pressure differential to the tool."}
{"text": "An article appeared on page 510, Volume 11, of the Philosophical Magazine in 1881 in which Alexander Graham Bell reported upon certain experiments which he conducted to study the sound-emitting properties of materials when exposed to the action of rapidly interrupted sunlight. Bell, among other things, introduced samples of material into a chamber and passed an intermittent beam of sunlight therein to produce audible effects. Bell observed that when the sample was in an acoustical resonant chamber the maximum response was achieved at a frequency of interruption equivalent to the resonant frequency of the chamber. While Bell was primarily interested in producing sound, he did recognize that this effect could be employed to study the properties of the material.\nIn the U.S. Pat. No. 3,948,345 which issued on Apr. 6, 1976, there is disclosed methods and means for analyzing solid and quasi-solid substances. This is accomplished by detecting signals which are established between a substance being tested and a surrounding fluid by a source of radiant energy whose amplitude is varied and whose frequency is scanned. An acoustic responsive transducer is located to sense the changes in the energy level of the surrounding fluid and transduces these acoustical energy changes into electrical signals that are characteristic of the substance being analyzed.\nIn the aforesaid U.S. patent there is disclosed in FIGS. 9-21 thereof the types of response graphs or curves obtained for a plurality of solid substances when analyzed by said patented system and method. Further reference to said patent is directed to obtain a more thorough understanding of said photoacoustic system for analyzing solid and quasi-solid substances.\nReference is also directed to the copending application, Ser. No. 567,640, filed Apr. 14, 1975 which is a continuation-in-part of U.S. Pat. No. 3,948,345, and which application discloses several embodiments of photoacoustic cells especially designed for use in said photoacoustic system of said U.S. Pat. No. 3,948,345.\nThe photoacoustic cell of this application is especially designed for use with the system and method of analysis of solid and quasi-solid materials as described in said aforesaid U.S. Pat. No. 3,948,345 and enables series of tests to be easily repetitively performed, said cell being relatively simple in structure and operation as compared to other sample test cells heretofore known.\nTherefore, an object of this invention is to provide a new and improved photoacoustic cell.\nAnother object of this invention is to provide a new and improved photoacoustic cell which is readily adaptable to perform repetitive tests on solid and quasi-solid substances.\nAnother object of this invention is to provide a new and improved photoacoustic cell which is applicable for a wide variety of uses.\nOther objects and advantages of the photoacoustic cell of this invention will become apparent to one skilled in the art upon reference to the following disclosure of a preferred embodiment."}
{"text": "Vascular endothelium is a cellular layer lining the inner part of blood vessels, including arteries and veins. Endothelium is presently looked upon as an important metabolically active autocrine/paracrine/endocrine organ that regulates cardiovascular function and maintains vascular homeostasis by: modulating vascular tone; regulating solute transport into cell components of the vessel wall, local cellular growth, and extracellular matrix deposition; protecting the vessel from the potentially injurious consequences of substances and cells circulating in blood; and regulating the hemostatic, inflammatory, and reparative responses to local injury. One of the main functions of endothelium is to produce or release substances, such as nitric oxide (NO), that control the behaviour of the blood vessels such as their dimensions, elasticity, permeability and reactivity, including the ability to constrict and dilate. Endothelium-derived mediators regulate not only blood flow and permeability vascular elasticity, reactivity and structure, but also local and systemic inflammatory response as well as thromboresistance of vessels. Vasoprotective endothelial mediators such as nitric oxide (NO), prostacyclin (PGI2) endothelium-derived hyperpolarising factor (EDHF), bradykinin (Bk) tissue plasminogen activator (t-PA), thrombomodulin (TM) or ADP-ase do exert antithrombotic, anti-inflammatory and vasoprotective action.\nOn the other hand, excessive production by endothelium of superoxide anions (O2-), isoprostanes, angiotensin II (ang II), endothelin 1 (ET-1), plasminogen activator inhibitor (PAI-1), tissue factor (TF), von Willebrandt factor (vWF), chemokines (e.g. monocyte chemotactive protein MCP-1), cytokines (e.g. IL-6), and increased expression of adhesion molecules (e.g. selectin P, ICAM-1) promote inflammation and thrombosis of vascular wall that may eventually lead to the development of atherosclerotic lesion. Accordingly, healthy endothelium is essential for undisturbed functioning of the cardiovascular system, while endothelial dysfunction leads to its various pathologies. In particular, endothelial dysfunction is pivotal to atherogenesis, it is present at the earliest stages (e.g. preceding angiographic or ultrasonic evidence of obstructive plaque) as well as later stages of arterial disease, contributing to clinical sequelae related to tissue damage (eg, ischemia, infarction, and organ failure).\nEndothelial dysfunction in most general terms refers to an impairment of the ability of the endothelial cell layer to produce an appropriate vasodilatory response to stimuli. Many studies provided evidence that endothelial dysfunction (assessed on the basis of the impairment of NO-dependent vasodilatation) may be regarded as prognostic factor for the development of adverse cardiovascular events. Indeed, relative risk for adverse outcomes is elevated approximately 10-fold when there is evidence of coronary or peripheral endothelial dysfunction.\nVarious conditions, including hypercholesterolemia, systemic hypertension, smoking, diabetes, congestive heart failure, pulmonary hypertension, estrogen deficiency, hyperhomocysteinemia, and the aging process itself, have been associated with impaired function (dysfunction) of endothelium. As a result, the vessel wall in these conditions may promote inflammation, oxidation of lipoproteins, smooth muscle proliferation, extracellular matrix deposition or lysis, accumulation of lipid-rich material, platelet activation, and thrombus formation. All of these consequences of endothelial dysfunction may contribute to development and clinical expression of atherosclerosis. The potential consequences of endothelial dysfunction further include coronary constriction or inadequate dilation during physical or mental stress, producing myocardial ischemia; plaque rupture and thrombosis, causing unstable angina or myocardial infarction; and reperfusion injury after thrombolysis.\nSeveral methods and apparatuses for non-invasive evaluation of the health of vascular endothelium in vivo have been developed.\nIn particular, methods are known that are based on monitoring the physiological conditions or characteristics of the arteries in the patient's limb after reactive hyperemia.\nReactive hyperemia is a physiological phenomenon that occurs in a patient after blocking (or occlusion) of a major artery. Such blocking or occlusion of artery in the limb, such as brachial artery, is typically done by inflating a blood pressure cuff slightly above systolic pressure for a period of about 5 minutes. Anoxia or severe hypoxia in the limb downstream from the occluded artery is usually a result of such blocking. Sudden release of the blocking causes endothelial cells to react by generating NO and dilating. The phenomenon of reactive hyperemia lasts up to 10 minutes before return to pre-test blood volume values. Blood flow is a characteristic of the artery, and under reactive hyperemia blood flow through an artery, vein or limb is significantly greater as compared with normal blood flow.\nCurrently the most popular method is flow mediated dilatation (FMD), a non-invasive technique based on monitoring of diameter of arteries after reactive hyperemia with a two-dimensional ultrasound and Doppler ultrasound. Its results correlate well with invasive coronary endothelial testing as well as with the presence and severity of coronary atherosclerosis. This technique is described for example in a review by S. Patel. And D. S. Celermajer, Pharmacological Reports 2006, 58, suppl. 3-7. However, this method is quite expensive, requires sophisticated equipment and highly specialized operators, is highly operator dependent and is poorly reproducible due to variability of measurements and poor resolution relative to arterial size. Hence, its use is limited and the method is not applicable on a more general basis.\nFor the purpose of assessment of vascular endothelial function changes of other physical parameters in response to reactive hyperemia have been also used, such as fingertip skin temperature (WO2005118516; N. Ahmadi et al. Int. J. Cardiovasc. Imaging (2009) 25:725-738), blood pressure in a finger (pulse wave amplitude) using plethysmography (EP1360929, EP1992282, WO00/57776, EP2110074), and peripheral arterial tone (WO2000/074551, WO2002/034105).\nNon-invasive technique for detection of endothelial dysfunction based on monitoring blood flow related changes in the level of a substance present in a limb after reactive hyperemia is disclosed in WO03/051193. The method involves blocking blood flow in the limb to stimulate endothelial function and then releasing the blood flow block to observe, measure and record said changes as a function of time, said changes being indicative of endothelial dysfunction. Said substance can be a tracer substance injected in a vein, such as a radiation emitter or a contrast agent, and the ingress of said tracer into the limb is detected and measured, for example by means of gamma ray detection. Tracer measurement in a pair of two laterally opposed limbs should be performed and the tracer presence compared between both limbs. Alternatively, a physical characteristics of the limb, such as temperature or color, or a property of a metabolic or other biochemical product circulating in the limb following the release of the blood flow block, such as O2, CO2 or reduced hemoglobin, is measured by a suitable technique. As suitable techniques there were suggested gas emissions across the skin surface within a cell placed on the skin surface, optical techniques, such as spectral analyzers or optical transmission/diffusion detectors, such as the visible-reflectance hyperspectral analysis, and EPR/NMR techniques. Either the appearance rate of a depleted substance or the disappearance (depletion) of an accumulated product can be detected. A rate of change of the measured parameter shortly after release of the occlusion or blockage is suggested as a primary factor in determining endothelial dysfunction. In the case of the use of a tracer, the rate for both the blocked limb and the contra-lateral control limb is measured.\nIt has been established that endothelial dysfunction is an early event and major risk factor for atherosclerosis and an important indicator for a medical professional, allowing for early diagnosis of the risk of cardiovascular disease.\nTesting endothelial function is therefore a highly desirable alternative for a diagnostic approach based on performing a set of various biochemical tests, especially in apparently healthy individuals, i.e. individuals not showing any signs of cardiovascular disease.\nA testing method is needed that would allow to evaluate function and detect any dysfunctions at an early stage of impairment in order to identify patients for prophylactic or therapeutic intervention to improve the dysfunction and/or for further more detailed and complicated diagnostic tests.\nThe need exists to provide a non-invasive test for evaluating endothelial function which would be reliable, easy to carry out and inexpensive, and thus applicable for tests in large patient populations, for example for screening purposes.\nThere is also a need for a simple, quick and non-expensive test that would allow to monitor and control the response of a patient to a medical treatment of cardiovascular disease."}
{"text": "Japanese Patent Application Publication No. 2008-238792 discloses a liquid ejecting device including a device body, and an ink cartridge. The device body includes a liquid ejection head and a subordinate tank. The ink cartridge has a liquid storage chamber and is adapted to be attached to and detached from the device body. The liquid storage chamber of the ink cartridge is provided with a sensor arm pivotally moved if the liquid level of the ink stored in the liquid storage chamber becomes equal to or lower than a predetermined level. The device body is provided with a sensor for detecting residual amount of ink. The sensor generates detection signals different from each other dependent on pivot position of the sensor arm. A controller of the liquid ejecting device determines residual amount of ink remaining in the liquid storage chamber of the ink cartridge on the basis of the detection signal outputted from the sensor. Further, annunciation is made to notify the user of the replacement of the ink cartridge with a new ink cartridge if the ink in the liquid storage chamber of the ink cartridge is consumed and the controller determines that the residual amount of ink is equal to or lower than the predetermined amount."}
{"text": "I. Field of the Invention\nThe present invention relates generally to a support structure that can be used to support shelving or other elements for carrying or supporting any desired item. More particularly, the present invention relates to a support assembly for use in, for example, a knock-down shelving system to adjustably support shelves.\nThe support assembly of the present invention ideally can be incorporated into a knock-down shelving system that includes a plurality of support posts for supporting one or more shelves at corner support assemblies thereof. The shelving system will include a sleeve or wedge member and a collar adapted to be secured to a generally cylindrical support post, that is one that is circular when viewed in radial cross-section. The sleeve has an inner surface that is configured to embrace the support post and an outer surface formed with a first axially outwardly inclined wedge portion and at least one second axially outwardly inclined wedge portion that, when viewed in such radial cross-section, extends at an angle to the first portion. Additionally, the collar is adapted to be secured to a member to be supported and to embrace the sleeve. The collar has a first wedge surface formed to mate with the first portion of the outer surface of the sleeve and at least one second wedge surface formed to mate with the second portion of the outer surface of the sleeve. Thus, when the sleeve embraces the post and the collar embraces the sleeve, axial loading of the collar in one direction causes the first wedge surface to mate with the first portion of the outer surface of the sleeve and the second wedge surface to mate with the second portion of the outer surface of the sleeve, thereby urging the sleeve toward the post in at least two generally radial directions.\nII. Description of the Prior Art\nShelving systems having adjustable height shelves and so-called “knock-down” type shelving systems are known, and each has utility in many applications. For example, a knock-down shelving system with adjustable height shelves may be used in the food service, industrial, commercial, hospital, and similar fields for storage of desired items.\nOne type of a well known knock-down shelving system is disclosed in U.S. Pat. Nos. 3,424,111 and 3,523,508 both to Maslow, and is manufactured and sold by InterMetro Industries Corporation, Wilkes-Barre, Pa., a company related to the assignee of the present invention. The shelving system disclosed in these patents has achieved great commercial success under the InterMetro trademark SUPER ERECTA SHELF®, and incorporates a plurality of cylindrical support posts each formed with a series of equally spaced, annular grooves on its outer surface. A basic shelving system might include four such posts to support one or more formed-wire shelves, with each shelf having a frusto-conically-shaped collar at each corner for receiving a support post. A two-piece interlocking sleeve fits around the support post. The sleeve features a rib on its interior surface for engaging one of the grooves on the support post and has a frusto-conically-shaped outer surface, which is widest at the bottom, designed to complement the shape of the shelf collars. The support posts fitted with sleeves are received in the collars of each shelf to assemble the shelving system. When assembled, the weight of the shelf and any items stored on it creates a radially-inwardly directed wedging force between the collars and sleeves, which brings the sleeves into tight contact with the posts.\nAnother type of commercially successful shelving system, sold and marketed under the trademark METROMAX® by InterMetro Industries, features a “knock-down” structure that incorporates triangular support posts. Such a system is the subject of U.S. Pat. Nos. 4,811,670; 4,946,350; 5,271,337; and 5,279,231.\nIn U.S. Pat. No. 4,811,670, a corner assembly for securing each corner of a shelf to the triangular support post includes a wedge member, a corner bracket structurally associated with the shelf, and a collar. The wedge member snap-fits onto the support post, and the collar and corner bracket form a sleeve that may fit around the wedge and support post, thereby to support the shelf by a wedging force.\nThe shelving systems shown in U.S. Pat. Nos. 4,964,350; 5,271,337; and 5,279,231 feature modular shelves in combination with the triangular support posts. The modular shelves include a rectangular shelf frame formed from two end beams connected to two side beams. A center beam may be inserted between the end beams, parallel to the side beams, to increase the load-bearing capacity of the system. A plurality of plastic shelf mats are adapted to be snap-fit onto the shelf frame. The shelf frame is secured to the support post by corner assemblies comprised of a corner portion of the end beam, a wedge member and a separate collar. A sleeve formed by the corner portion and the collar is seated on the support post and wedge member and secured by a wedging action. Two lock cylinders lock the collar to the corner portion to secure the sleeve.\nStill another type of commercially successful shelving system manufactured and sold by InterMetro Industries is disclosed in U.S. Pat. No. 6,113,042 (Welsch). In one embodiment, this system includes a two-piece wedge assembly configured to embrace a circular support post. The wedge assembly is formed of a sleeve and a wedge that are snap-fit or otherwise joined together about the support post. The wedge has a single planar wedge surface that tapers outwardly from its upper end to its lower end, such that the lower end is wider, and extends toward the interior of the shelving system. A collar, having a rear section that outwardly inclines from top to bottom to complement the slope of the single wedge surface, is seated on the wedge assembly and secured by wedging action.\n(Each of the patents mentioned above is incorporated in its entirety herein by reference.)\nDespite the significant utility and commercial success of the above-described shelving systems, a need exists for an even more highly stable shelving system, utilizing a cylindrical support post, that is easily assembled and has shelves that are easily adjusted to different heights without the need for special tools, and in which the shelves are secured in a static manner to provide a load carrying capacity suitable for heavy-duty use. In particular, it is desirable to provide a system in which a support member can be secured to a cylindrical support post so as to provide stability and rigidity in multiple radial directions relative to the axis of the post."}
{"text": "The present invention relates to a Voltage Source Converter having at least one phase leg connecting to opposite poles of a direct voltage side of the converter and comprising a series connection of switching cells, each said switching cell having on one hand at least two semiconductor assemblies connected in series and having each a semiconductor device of turn-off type and a free wheeling diode connected in parallel therewith and on the other at least one energy storing capacitor as well as two terminals connecting the cell to adjacent cells in said series connection of switching cells, a mid point of said series connection of switching cells forming a phase output being configured to be connected to an alternating voltage side of the converter, each said switching cell being configured to obtain two switching states by control of said semiconductor devices of each switching cell, namely a first switching state and a second switching state, in which the voltage across said at least one energy storing capacitor and a zero voltage, respectively, is supplied across said two terminals of the switching cell, for obtaining a determined alternating voltage on said phase output, in which said series connection of switching cells includes inductance means.\nSuch converters with any number of said phase legs are comprised, but they have normally three such phase legs for having a three phase alternating voltage on the alternating voltage side thereof.\nA Voltage Source Converter of this type may be used in all kinds of situations, in which direct voltage is to be converted into alternating voltage and conversely, in which examples of such uses are in stations of HVDC-plants (High Voltage Direct Current), in which direct voltage is normally converted into a three-phase alternating voltage or conversely, or in so-called back-to-back stations in which alternating voltage is firstly converted into direct voltage and this is then converted into alternating voltage, as well as in SVCs (Static Var Compensator), in which the direct voltage side consists of capacitors hanging freely. However, the present invention is not restricted to these applications, but other applications are also conceivable, such as in different types of drive systems for machines, vehicles etc.\nA Voltage Source Converter of this type is known through for example DE 101 03 031 A1 and WO 2007/023064 A1 and is as disclosed there normally called a multi-cell converter or M2LC. Reference is made to these publications for the functioning of a converter of this type. Said switching cells of the converter may have other appearances than those shown in said publications, and it is for instance possible that each switching cell has more than one said energy storing capacitor, as long as it is possible to control the switching cell to be switched between the two states mentioned in the introduction.\nAnother Voltage Source Converter of this type is known through U.S. Pat. No. 5,642,275 used in a Static Var Compensator, in which the switching cells have a different appearance in the form of so-called full bridges.\nThe present invention is primarily, but not exclusively, directed to such Voltage Source Converters configured to transmit high powers, and the case of transmitting high powers will for this reason mainly be discussed hereinafter for illuminating but not in any way restricting the invention thereto. When such a Voltage Source Converter is used to transmit high powers this also means that high voltages are handled, and the voltage of the direct voltage side of the converter is determined by the voltages across said energy storing capacitors of the switching cells. This means that a comparatively high number of such switching cells are to be connected in series for a high number of semiconductor devices, i.e. said semiconductor assemblies, are to be connected in series in each said switching cell, and a Voltage Source Converter of this type is particularly interesting when the number of the switching cells in said phase leg is comparatively high. A high number of such switching cells connected in series means that it will be possible to control these switching cells to change between said first and second switching state and by that already at said phase output obtain an alternating voltage being very close to a sinusoidal voltage. This may be obtained already by means of substantially lower switching frequencies than typically used in known Voltage Source Converters of the type shown in FIG. 1 in DE 101 03 031 A1 having switching cells with at least one semiconductor device of turn-off type and at least one free-wheeling diode connected in anti-parallel therewith. This makes it possible to obtain substantially lower losses and also considerably reduces problems of filtering and harmonic currents and radio interferences, so that equipment therefore may be less costly.\nAccordingly, the invention is directed to converters with so-called cascaded half or full bridges.\nIt is an ongoing attempt to improve such converters already known for instance by reducing the complexity and/or the size thereof."}
{"text": "1. Field of the Invention\nThis invention relates to a method and apparatus for providing a single service destination in a telecommunications network. In particular, but not exclusively, the invention relates to the provision of a single mailbox for multi-service users having access to a plurality of different telephony services.\n2. Description of the Related Technology\nKnown voicemail systems allow a calling party to leave a message for a called party when the called party is unable or unwilling to answer the call. If a call is not answered by the called party, the calling party is directed to a local messaging system or a remote messaging centre where the calling party may record a message for subsequent retrieval by the called party. Such systems are known in both fixed and mobile telephone networks and also in data networks such as the Internet.\nFIG. 1 shows a prior art telecommunications network. The telecommunications network includes a Public Switched Telephone Network (PSTN) 102 to which a telephone 100 (telephone A) is connected and which is connected to a service provider network 104. Service provider network 104 includes a local telephone exchange switch 110 which serves a telephone 106 (telephone B). The user of telephone B has a voice messaging service provided by the service provider who operates service provider network 104. The voice messaging service is implemented via a messaging centre 108 (MCB) which is connected to local telephone exchange switch 110. Operation of messaging centre 108 is under the control of the service provider of network 104.\nFIG. 2 is a flow diagram according to the prior art which illustrates the steps involved in a user of telephone A leaving a voice message for the user of telephone B, for an arrangement in which the user of telephone B has a voice messaging service provided by the service provider who operates network 104 (shown by annotation 200 in FIG. 2).\nWhen telephone A is used to make a call to telephone B, the telephone dialing number for telephone B is dialed. This initiates a call that is routed via PSTN 102 to local telephone exchange switch 110 which serves telephone B located in service provider network 104. The call initiation results in a signaling message (shown by step 2a) traversing PSTN 102 which contains a calling line identifier (CLI) (also referred to as a caller ID) for telephone A (denoted CLI: A) and destination dialing number for telephone B (denoted DN: B). When local telephone exchange switch 110 receives the signaling message of step 2a, it rings telephone B as shown by step 2b. \nIn this case, the user of telephone B has subscribed to CLI services and the user of telephone A has not withheld their CLI when making the call, so the signaling message of step 2b contains a CLI for telephone A.\nIf telephone B is not answered and the user of telephone A wishes to leave a voicemail message for the user of telephone B, then after a pre-configured delay, local telephone exchange switch 110 forwards the call to messaging centre 108. The signaling message which is sent to messaging centre 108 (shown by step 2c) contains the originating CLI for the call (denoted CLI: A) and a destination dialing number for messaging centre B (denoted DN: MCB). The signaling message of step 2c also contains a redirecting number, which in this case is the number of the line forwarding the call to messaging centre 108. The redirecting number allows messaging centre 108 to identify in which mailbox to deposit the voicemail message from the user of telephone A: in this case telephone B's mailbox.\nFIG. 3 is a flow diagram according to the prior art which illustrates a user of telephone B retrieving a voice message from their mailbox using telephone B. The user of telephone B has a voice messaging service provided by the service provider who operates service provider network 104 (shown by annotation 300 in FIG. 3).\nTo retrieve their voicemail messages, the user of telephone B calls messaging centre 108, shown by step 3a. This may be done by dialing a short code associated with the messaging centre. A short code may be a three digit short code such as 121, 123 or 901 on UK mobile networks, whereas in US mobile networks, a user may dial their own mobile number from their mobile telephone, which is then forwarded to the messaging centre as their mobile telephone is engaged. Alternatively, the user can dial the full telephone dialing number of messaging centre 108.\nThe signaling message for step 3a contains a CLI for telephone B (denoted CLI: B) and a destination dialing number for the messaging centre (denoted DN: MCB). Messaging centre 108 checks the CLI of the incoming call, and routes the caller to the appropriate mailbox, i.e. that of the user of telephone B. The user can then either listen to their messages, or, if the system is so configured, may have to enter a Personal Identification Number (PIN) first to prevent unauthorised access to their mailbox.\nFIG. 4 is a flow diagram according to the prior art which illustrates a user of telephone B retrieving a voice message from their mailbox. The user of telephone B has a voice messaging service provided by the service provider who operates service provider network 104 (shown by annotation 402 in FIG. 4). In this case, the user retrieves their voice message from a telephone D (see annotation 400) which is different to telephone B.\nTo retrieve their voicemail messages, the user of telephone B calls messaging centre 108 from telephone D. To retrieve voicemail messages from a telephone different to the telephone B, the user dials the full telephone number of messaging centre 108 that hosts their mailbox which causes a signaling message to be sent to messaging centre 108 in step 4a. \nThe signaling message for step 4a contains a CLI for telephone D (denoted CLI: D) and a destination dialing number for the messaging centre (denoted DN: MCB). Messaging centre 108 checks the CLI of the incoming call, but does not recognize it as it identifies a telephone foreign to their service, i.e. telephone D. Messaging centre 108 therefore cannot connect the user directly to their voice mailbox. Instead, messaging centre 108 prompts the user to enter the number of the phone which has a mailbox on that messaging centre, i.e. that of telephone B, and then asks for the user's PIN to authorise access.\nThe above prior art systems which allow voicemail to be left in mailboxes by calling parties are adequate when a user has a single telephony service. However, if a user has more than one telephone service, for example a mobile telephony service and a landline telephony service, then the user has to deal with different mailbox services hosted by each of the different service providers. This can be impractical for the user as they have to coordinate the messages left in their different mailboxes. Further, the user has to deal with different access methods, different access codes and different user interfaces for the different mailbox services, which can lead to an inefficient use of the user's time and may lead to increased costs if each of the mailbox services must be paid for individually.\nThe above prior art methods for retrieving voicemail from a mailbox have a problem in that, if a user wishes to retrieve mail from their mailbox using a telephone service not provided by the same provider as who provides their mailbox service, then it is either not possible for the user to gain such access, or the user must enter a full telephone dialing number associated with the relevant mailbox service. This can be time consuming for the user, and especially inconvenient if the appropriate number is not readily to hand.\nU.S. Pat. No. 7,023,967 relates to a system for sharing voice messaging support which includes a voice messaging system coupled to a wireless network and a wireline network. A mailbox can be shared between telephones on each of the different networks. In order to implement this prior art system, each of the wireline and wireless telephony services must be configured to incorporate the shared mailbox functionality.\nIt would therefore be desirable to provide improved mailbox services for multi-service users who have access to a plurality of telephony services. A single mailbox for the multiple telephony services would be beneficial, since it would facilitate efficient access methods, both for calling parties leaving messages in the mailbox and for subsequent retrieval of the messages by the multi-service users. Further, it would be desirable to provide such mailbox services with a minimum of reconfiguration and disruption to the existing telephony services and telephony devices."}
{"text": "Projectiles currently used for the M203/M320 gun launcher devices include the M433 High Explosive Dual Purpose Round, M406 High Explosive Round, M583A1 Star Parachute Round, M585 White Star Cluster Round, M713 Ground Marker Round, M781 Practice Round, M651 CS Round and M576 Buckshot Round. The M203 is not a stand-alone gun. It must be attached to an M16 or M4 rifle, for instance, and fires a family of 40 mm low velocity projectile rounds. In order to load a 40 mm round into an M203, the forward sleeve of the gun tube is slid forward, and a 40 mm round is pushed into the sleeve. The sleeve is then slid backwards and locked into position, ready to fire. The barrel of the M203 is rifled, which spins up the projectile as it travels down the barrel tube, and the projectile leaves the gun tube with a high spin rate.\nThe 40 mm cartridge case is typically not reused or recycled after firing. If the cartridge case becomes stuck in the gun tube, there is a tool to help the war-fighter push it from the tube. The M320 gun launcher device fires the same 40 mm low velocity ammunition as does an M203. There are several improvements that the M320 has compared to the M203. The M320 has stand alone single shot capability and may be fired by a war-fighter without attaching to a gun. It also has the capability to be attached to and fired from an M16/M4 rifle, similarly to an M203 gun launcher device. A major improvement of the M320 is the ability of the firing tube to open sideways. By opening sideways, longer ammunition can be loaded into the M320 gun tube as compared to loading into an M203. While the approximate maximum length of the projectile is approximately five inches in the M203, projectiles several inches longer can be loaded into an M320. Both M203 and M320 have rifled barrels which induce spinning in the projectile as it travels through the gun tube.\nThere exists a great need to develop a new family of 40 mm projectiles that can have extended range through increased velocity and projectile shape to over 1000 meters; can provide capability for an add-on optional guidance-navigation and control (“GN&C”) subsystem; can increase lethality through enhanced precision and accurate strikes; can increase projectile length which can accommodate more fragments and an increased amount of explosives and/or energetics; and can provide a novel system with much longer rear end fins, that can allow for greater stability, in applications which will no longer need any appreciable, or no, projectile spin for flight accuracy.\nCurrently, low velocity (“LV”) 40 mm projectiles are length restricted by the presence of the loading system that is in the M203 gun. The 40 mm LV projectiles do not fly as well as desired, due to the inability to add fins for providing stability. Instead, 40mm LV projectiles are spin stabilized, which both do not fly as well as desired and also can not employ guidance systems due to these projectiles' extensive pitch, yaw and drag. Current guidance-navigation and control (“GN&C”) subsystems can't function in a spin environment. Guided ammunition would need to have no spin (or very little spin) so that the seeker or camera therein, e.g., can see the target clearly while in flight. However, the conventional M203 or M320 grenade launcher guns do have spin since they have internal rifling intended to spin up the projectiles (at approximately 1800 revolutions per second) for stabilization of flight. Clearly then, a way must be found to avoid the spin which is inherent in conventional M203 or M320 grenade launcher guns when optional guidance is desired to be added to the projectile.\nRecently developed new 40 mm projectiles of several varieties provide better performance, extended range, with reduced or no spin at all, and have the ability to have added on guidance and control systems to hit stationary targets. These novel technologies all allow for better performance, yet while reducing pitch, yaw, and/or drag. The stability of these projectiles are still somewhat limited by the length and area of their fins. These projectiles will fly well and have the ability to add guidance and control with small canards added in the front of the projectile, in a canard actuating system (“CAS”). Smaller (width and/or length) canards however would mean less control authority or less of an ability to make large course corrections while in flight. The projectiles' center of gravity (“CG”) and center of pressure (“CP”) needs to be a certain distance apart in order to provide a needed stability margin. In the case of fin stabilized projectiles, the CG needs to be in the front of the projectile while the CP needs to be in the rear of the projectile. Larger canards, and changing the position of the canards away from the front of the projectile, move the CP forward and begin to compromise or reduce the stability margin. By incorporating longer fins on a projectile with more surface area however, the CP is moved back and the projectile gains a better stability margin. With this greatly improved stability margin, larger canards can be positioned where needed. This improvement now allows for a projectile to have larger course corrections and hit harder targets such as moving targets. Optional guidance and control can now be added to the front or side of a projectile utilized for battle. The optional guidance and control provides increases accuracy and precision hits on targets up to and beyond 1000 meters (a goal which is sought). The obturator on the projectile prevents gases from blowing past it; it seals in the gases allowing a projectile to move with a constant rate push as it travels down the barrel of a gun launcher. The obturator must be positioned immediately behind of the cartridge case. This is because the M320 and M203 have both a section to seat the cartridge case, and then also enough room still, for the obturator before the rifling part of the gun tube is encountered. The obturator is higher than (near the end of) the projectile and therefore is found at the top of the gun tube. If the obturator were made too long in dimension, then the cartridge could not be loaded. This is since the obturator cannot be pushed into the rifling just by hand. Currently, there is an approximately one inch long of space in between the back of the obturator and the front of the (inside of the) cartridge case. This can only allow for fins that are one inch long therefore. The projectile continues on past the obturator by a distance of approximately 1/16 inch (press fit into the cartridge case) and then slopes down at an angle of 15 degrees, toward the end of the projectile. There is currently provided a housing for the fins at the end of the projectile, which housing functionally attaches the fins to the projectile. The fins physically now extend from the housing and up until the beginning of the interference fit with the cartridge case. Therefore the fins are approximately one inch long. The fins open up after exit from the gun. These finned projectiles have tested with excellent flight results, however, to add in a guidance and control (G&C) canard system as desired, the fins would not allow for larger and more aggressive canards and positioning. These projectiles are good for stationary targets. To hit moving and stationary targets, fins that are two inches or longer are needed. Therefore, what is greatly needed is to provide a low velocity 40 mm projectile with two or more inch long fins to allow for aggressive guidance and control CAS systems. In addition, aeroballistic modeling shows that the rear end of the projectile will have less drag with slopes on the projectile that are less than eight degrees, rather than the currently mentioned fifteen degrees, and this type of projectile is also greatly needed."}
{"text": "When taking pictures a user of an image capture device such as a camera for example may or may not put care into framing a subject of the image accurately. Hence, in some situations pictures are taken quickly and are generally referred to as ‘snaps’ or ‘snapshots’, whilst in others a long time can be spent carefully framing a subject in order to obtain a desired image composition.\nDifferences in picture taking behaviour can be the result of a combination of factors. For example, some camera users always take pictures very quickly, perhaps due to lack of training or through not wishing to spend a long time “behind the camera”, whilst other users take time and trouble when the situation permits or a good result is critical, but on other occasions are forced to take a quick snap.\nA further distinction in picture taking behaviour relates to the subject itself. In the case of moving objects, a photographer may wait for a desired moment to capture an image. Some examples include: a subject scoring a goal, eye-contact with the user, and the sun emerging from behind a cloud and so forth. In these situations, attention is more focused on capturing a desired moment in time, rather than ensuring that a subject of the captured image is well framed.\nCurrent cameras do not recognize the picture taking behaviour of a user, other than by the direct operation of the camera controls. As such, the image processing functionality of the camera is generally determined only by its settings and the content of the captured image. This precludes the possibility of making image processing responsive to the way the camera is being used.\nA particular problem with this rigid approach to image processing is that processing steps such as autocropping are either ‘on’ or ‘off’, as generally set by user choice —the term “autocropping” refers to algorithms which automatically generate one or more crops of a photographic image based on an algorithmic assessment of likely areas of interest, or salient areas, within a captured image. The autocrop location can be selected to include the areas of interest, remove distractions or large boring areas, and produce a pleasing composition for example.\nA number of autocropping approaches are known. U.S. patent application Ser. No. U.S. Pat. No. 5,978,519 is an example. It assumes that the “background” region of a captured image has relatively low intensity variance compared with the foreground. A threshold is selected based on the variance of blocks, and a bounding box of the subject is then used for cropping.\nA more recent approach may be found in United Kingdom Patent Application No. GB0031423.7 which uses a saliency map determined by region segmentation and colour ‘unusual-ness’ to identify both subject areas and distraction areas. Multiple crops are suggested based on alternative combinations of subjects, in each case excluding distractions and selecting locally optimal crop boundaries.\nEuropean Patent Application No. EP1158464A1 also uses a saliency map derived from region segmentation.\nSome of the above algorithms can have hints or constraints applied by an external source (typically a person interacting with the algorithm). These provide the algorithm with guidance to ensure that some areas are always preserved, but do not provide an automatic determination of when autocropping may be appropriate."}
{"text": "The present invention relates generally to the field of graphical user interface (GUI) development, and more particularly to computer-automated testing of GUI mirroring.\nGUIs support a variety of natural languages in order to target users from around the globe. Supporting different languages requires that a GUI display a language in its normal direction flow. For example, English and French are read from left to right (LTR), while Arabic, Hebrew, and Persian are read from right to left (RTL). User convenience dictates that the layout of a GUI for an RTL language be mirrored with respect to the corresponding LTR layout. GUI mirroring generally involves reflecting the position of GUI elements about a vertical axis passing through the center of the screen. Nested GUI elements will be similarly reflected relative to their containing elements. In addition, the direction flow of text elements may be switched, based on whether the language is LTR or RTL.\nGUI elements are those elements used by GUIs to visually represent information stored in a computer. These generally facilitate the use of computer software. Examples of GUI elements are windows, for example, folder windows and message windows; and control elements, or widgets, which facilitate specific user-computer interactions, such as buttons, menus, and text boxes.\nTo ensure that mirroring is implemented correctly, such implementations are generally tested manually. A human tester may preview a mirrored GUI to check the correct placement of GUI elements and report if any are incorrect. Such manual efforts, which need to be redone each time a change to the GUI design occurs, tend to be time consuming and subjective—small deviations from the correct positions may be easily missed by the naked eye. Moreover, certain device characteristics may cause a GUI to appear different when mirrored on different display devices."}
{"text": "The present invention relates to a nonvolatile semiconductor storage device that has a means for executing erase in blocks and a storage contents erase method therefor.\nThere has recently been a growing demand for a nonvolatile semiconductor storage device capable of executing erase in blocks, intended mainly for portable equipment.\nThe block erase of a conventional nonvolatile semiconductor storage device capable of executing erase in blocks will be described below with reference to FIGS. 9A, 9B, 10A, 10B, 11, 12 and 13.\nFIG. 9A shows a typical memory cell structure of an EEPROM (electrically erasable programmable ROM), which is a nonvolatile semiconductor storage device. As shown in this figure, the memory cell has a MOS transistor structure of a two-layer gate constructed of a control gate 701 and a floating gate 702. This memory cell is covered with an insulating film of SiO2 or the like. This insulating film serves to provide electrical insulation of each portion that constitutes the memory cell, a function as a capacitor and protection from the external environment. A state in which a comparatively large number of electrons are injected into the floating gate 702 serves as a written state, where the memory cell has a high threshold value. A state in which a comparatively small number of electrons are injected into the floating gate 702 serves as an erased state, where the memory cell has a low threshold value. This difference between the threshold values is utilized for information storage. An operation for determining whether each memory cell is in the written state or in the erased state is a read operation.\nWrite into the memory cell, i.e., a transition from the erased state to the written state is achieved by injecting electrons into the floating gate.\nThere have been put into practice several methods for this purpose, and the write by injecting channel hot electrons (CHE) described below is most general. In concrete, by applying a high voltage (10 V, for example) to the control gate 701, applying a high voltage (6 V, for example) to the drain 705 and setting the source 703 to 0 V, a channel is formed to flow a large current between the drain 705 and the source 703. That is, electrons move from the source 703 to the drain 705. The electrons that have moved from the source 703 to the drain 705 become electrons in a high energy state due to the high voltage of the drain 705. If the energy at this time exceeds the energy barrier of an insulating film 704, then the electrons can move to the floating gate 702. With this mechanism, the memory cell is put in the written state by the injection of electrons into the floating gate 702.\nOn the other hand, the erase of the memory cell can be achieved by extracting the electrons accumulated in the floating gate 702. Several methods have been put into practice, and a method for extracting electrons from the source 703 is most general. According to this method, the electrons are moved by the tunnel effect from the floating gate 702 to the source 703, i.e., the memory cell is erased by, for example, making the control gate 701 have a voltage of 0 V, making the source 703 have a high voltage (12 V, for example) and floating the drain 705. This is the erase method called the high voltage source erase.\nThere is also put into practice a negative voltage gate erase, which is a kind of source erase and able to suppress the source voltage low. According to this erase method, the potential of the floating gate 702 is lowered by applying a negative voltage (xe2x88x9210 V, for example) to the control gate 701, applying a high voltage (5 V, for example) to the source 703 and floating the drain 705. According to this erase method, a similar tunnel effect can be obtained by a source voltage lower than that of the method of making the control gate 701 have a voltage of 0 V, and the memory cell erase can be achieved. This negative voltage gate erase is also an erase method that applies a high voltage to the source. In order to distinguish this method from the source erase that is not the negative voltage gate erase, there will be provided the description that the former is referred to as a negative voltage gate erase and the latter is referred to as a high voltage source erase in distinction.\nNext, FIG. 9B shows the array structure of a flash memory. This figure shows the array structure of a NOR type flash memory, which is a typical flash memory. Row select lines 711, 712, 713, 714, 715, 716, . . . are connected to the control gates of a plurality of memory cells, and column select lines 732, 731, . . . are connected to the drains of the plurality of memory cells. The plurality of row select lines 711, . . . , 716, . . . and the plurality of column select lines 732, 731, . . . constitute a matrix, forming a memory array. In the flash memory, the memory cells in an identical block share a source line 741, and this facilitates easy batch erase of the cells in a block and also enables a substantial reduction in the memory array area.\nDuring write in this flash memory cell array, the write is effected only on the cell of which both the row select line and the column select line are selected, and therefore, write in bits can be achieved. In the case of division into blocks that respectively have a common source line as in a flash memory, erase is executed by block erase for erasing in a batch all the memory cells in a block. It is to be noted that the block erase can be achieved with a lower source voltage than in the case of the high voltage source erase when the negative voltage gate erase is used, and therefore, only the memory cell to the gate of which a negative voltage is applied can be erased. It is also possible to execute sector erase for selectively erasing only a specified sector by dividing the block that shares a source line into sectors that are more minute erase units.\nDuring the source erase of an EEPROM, a tunnel current between bands (BTBT: Band-To-Band Tunneling current, referred to as a BTBT current hereinafter) described in detail later cannot be avoided in the overlap region of the floating gate and the source diffusion layer. Accordingly, there is a reduction in the current efficiency of the erase operation, i.e., the ratio of electric charges to be extracted from the floating gate with respect to the charges consumed by the high voltage applied to the source.\nThis BTBT current will be described with reference to FIGS. 10A and 10B.\nFIG. 10A schematically shows the state in the vicinity of the source 703 in the aforementioned erase operation. This erase operation is achieved by extracting electrons from the floating gate 702 to the source 703 by the FN (Fowler-Nordheim) tunneling phenomenon. An electron A inside the floating gate 702 is moved to the source 703 by the voltage applied during erase according to the FN tunneling phenomenon. This movement of the electron A is the erase operation.\nHowever, by applying a high voltage to the source 703, a potential slope, i.e., band bending is caused by an electric field that concentrates on the surface and its vicinities of the overlap portion of the floating gate 702 and the source 703. This increases the electron potential in a valence band, and if the potential becomes higher than the conduction band of the N+ region, then there appears an electron (electron B in FIGS. 10A and 10B) that passes through the bandgap between the valence band and the conduction band by the band-to-band tunnel effect and moves to the conduction band. This is the BTBT current, occurring concurrently with a hole. FIG. 10B shows an energy band diagram of the state of electron energy at this time. The electron A is an electron that moves from the floating gate 702 to the source 703 during erase. The electron B moves in the direction of arrow as a consequence of an increase in the electron potential in the valence band due to band bending \"psgr\"s. That is, the BTBT current is generated. Due to this BTBT current, the current efficiency of the erase operation of source erase deteriorates. The BTBT current not only reduces the current efficiency of the erase operation but also generates a hole, becoming an important factor to deteriorate the reliability of the memory cell. That is, due to the trap of the hole, which has been generated by the band-to-band tunnel effect, in the insulating film, the reliability of the insulating film 704 is degenerated.\nSince the sources of all the memory cells in a block are connected to an identical source line in the flash memory in which the block batch erase is executed, an increase in the peak current due to the BTBT current becomes a serious problem. The reason for the above is that a portion that belongs to the consumption current during the erase operation and depends on the BTBT current is roughly proportional to the number of memory cells to which the erase voltage is simultaneously applied to the source. In addition, with regard to the flash memory, due to the fact that the device operable on a single power supply voltage is mainly used for the sake of improvement in operability and the fact that the power supply employed in portable equipment on the big market of the flash memory is limited, it is often the case where a voltage boosted from the power supply voltage by a charge pump circuit is used as a source voltage to be used for source erase. In order to generate a specified high voltage regardless of the low power supply voltage, a larger current is consumed. In addition, the boosting efficiency of a voltage booster circuit, i.e., a ratio of an output electric power to an input consumption power is generally reduced in accordance with a reduction in the power supply voltage. Therefore, an increase in the consumption current due to the BTBT current becomes still more serious problem in accordance with the adoption of a single power supply voltage and a low voltage.\nNext, FIG. 11 shows an example of a booster circuit for generating a high voltage to be applied to the source during erase. In this FIG. 11, only the part necessary for explaining the boost is described.\nThe booster circuit shown in FIG. 11 includes a charge pump circuit 808 for effecting boost from a power supply voltage 801 and generating a boosted voltage 899. In this pump circuit 808, a plurality of sets of a capacitor and a transistor such as the sets of a capacitor 821 and a transistor 822 are connected in series in a plurality of stages. Then, a set constructed of a capacitor N3 and a transistor N4 serves as the final stage of the sets connected in series.\nThe oscillation circuit 802 of FIG. 11 alternately puts two signal lines 803 and 804 into an enable state. A buffering circuit 805 alternately thrusts nodes 806 and 807 up to the power supply voltage on the basis of signals from these signal lines 803 and 804. For example, in the capacitor 821, the node 806 being thrust up raises the potential of the other electrode, and the boosted charges pass through the transistor 822. This charges, which have passed through the transistor, are further boosted by the operation of the capacitor 823. By repeating this operation, the booster circuit shown in FIG. 11 can generate the higher voltage 899 on the basis of the power supply voltage 801.\nThen, since the transistors 822, 824, . . . , N2, N4 and so on are connected so as to prevent the backward flow, the charges of which the potential has been raised by the capacitors flow toward an output node 888. By repeating this operation, the booster circuit of FIG. 11 generates the high voltage. The fundamental operation of the booster circuit is described above.\nNext, FIG. 13 shows one example of the current waveform when the erase voltage is applied. It is needless to say that the consumption current is not entirely attributed to the BTBT current. The contents of the consumption current include a variety of elements such as a current consumed by a write state machine for the control of the erase operation and a current used for the voltage control of other than the source. However, during source erase, the consumption current attributed directly to the BTBT current is especially large, and it can be safely said that the greater part of the current that varies on the time base is attributed to the BTBT current. Therefore, no explanation is provided hereinafter for the current attributed to other than the BTBT current so long as the need does not arise.\nThe current efficiency of the erase operation is very low due to the BTBT current. Moreover, the potential of the floating gate rises as the erase progresses. This reduces the electron potential of this region, suppresses the peak of the potential in the valence band on the surface of the N+ region and reduces the BTBT current. Therefore, the peak of the consumption current is located immediately after the application of the consumption voltage during the source erase, and the consumption current reduces according to the progress of the erase. Generally, in order to control the reduction in the power supply voltage attributed to the consumption of current, a capacitor is connected to the power supply outside the device. However, the erase time of the flash memory is very long (several milliseconds to several seconds), and it is not easy to secure a capacity sufficient for handling a large current for a long time. Due to this large peak current, it has been difficult to use a flash memory in a system such as portable equipment, which has only a small current supply capacity.\nA concrete example of the execution of the negative voltage gate erase will be described more in detail with reference to FIG. 12. In the following description, a group constructed of a single or a plurality of row select lines is referred to as a row select line group, and a group constructed of a single or a plurality of column select lines is referred to as a column select line group. FIG. 12 shows only the part necessary for explaining the block erase of a nonvolatile semiconductor storage device capable of executing block erase. There are shown a memory array 910 that becomes an erase unit, a row decoder 920 for controlling the row select line group 921, a negative voltage 901 supplied to the row decoder, a debooster circuit 960 for generating the negative voltage 901, a column decoder 930 for controlling the column select line group 931, a source control circuit 940 for controlling the source line 941, a high voltage 902 supplied to the source control circuit 940 and a booster circuit 950 for generating the high voltage 902. In this case, the booster circuit 950 and the debooster circuit 960 have a construction similar to that of the aforementioned booster circuit shown in the FIG. 11. The row select line group 921 shown in FIG. 12 includes all of the row select lines connected to the memory cells included in the memory array 910. The column select line group 931 includes all of the column select lines connected to the memory cells included in the memory array 910.\nWhen the erase of the memory array 910 is executed, the row decoder 920 applies the negative voltage 901 generated by the debooster circuit 960 to all the row select lines included in the row select line group 921. The source control circuit 940 applies the supplied high voltage 902 to the source line 941 of the memory array 910. The column select line group 931 is controlled so as to be floating by the operation of the column decoder 930 or a circuit that controls the drain voltage and is connected to the column decoder 930. Control of the drain voltage, which is not especially important in terms of the construction of the present invention, is not shown in FIG. 12. By the application of these voltages, all the memory cells included in the memory array 910 receive the negative voltage applied to their control gates and the high voltage applied to their sources and have the drain floating. This is the conventional erase operation, and the consumption current roughly becomes as shown in FIG. 13.\nDuring the negative voltage gate erase, the sector erase can be achieved as described hereinabove. According to this method, erase is executed not on the entire block, and a negative voltage is applied only to a part of the row select lines connected to the erase block. The memory cells connected to the row select line to which this negative voltage is applied are erased, and a BTBT current is generated. On the other hand, the memory cells connected to the row select line to which the negative voltage is not applied are not erased since the electron potential of the floating gate is lowered. Furthermore, since the electron potential of the valence band in the vicinity of the N+ surface of the source is not significantly raised, the BTBT current is largely reduced. However, according to this method, the erase time is increased for the execution of the erase of the entire block. For example, when the erase is executed half-and half with regard to the number of sectors included in a block in order to suppress about half the peak current, twice the erase time is simply required.\nMoreover, the BTBT current can be reduced by the channel erase method. According to the channel erase, the erase operation is achieved by extracting the electron in the floating gate to the substrate with a negative voltage applied to the control gate and with a positive voltage applied to the substrate. According to this method, the BTBT current can almost be removed. However, due to the necessity of substrate potential control, there is needed a substantial change of the manufacturing process and the control method, which cannot easily be achieved.\nMoreover, the BTBT current can also be suppressed by setting low the source voltage during erase. However, there is a demerit that the time required for the erase largely increases by simply lowering the source voltage.\nAccordingly, a technology for restraining the increase of the erase time to a minimum while suppressing the peak current is soft erase. This soft erase, which is originally a technology intended for improving the reliability of the memory cell, also serves as a means for suppressing the BTBT current ascribed to the hole generating mechanism during erase and is also effective for a reduction in the erase current.\nDuring this soft erase, the source voltage is set lower than usual for a certain period from the start of the erase operation, and the source voltage is set back to the normal voltage after the progress of the erase operation to a certain extent. Since the source voltage is low immediately after the application of the erase voltage when the BTBT current is maximized, the peak current is suppressed. When the source voltage is set back to the normal erase voltage, the peak current does not increase so much because the erase has progressed to a certain extent.\nHowever, it is very difficult to quantitatively grasp the correlation between the source voltage and the peak current value if the variations of the process, the deterioration of the memory cells and so on are taken into consideration. Moreover, the same thing can be said for the quantitative correlation between the erase time and the source voltage. Therefore, it can be qualitatively estimated to reduce the peak current by the soft erase, but a quantitative determination is difficult.\nAccordingly, as a means for enabling the quantitative estimation of the suppression of the consumption current peak, there is proposed a method for limiting the consumption current of the source voltage generation circuit. This method is a method for providing the source voltage generation circuit with a current limiting means. This technology pays attention to the fact that the consumption current attributed to the BTBT current is consumed by the booster circuit for generating the source voltage. The high voltage generated by the booster circuit flows a current to the substrate in the form of the BTBT current. However, when the electric charges lost by the BTBT current cannot be replenished, the source voltage is lowered due to limitations on the consumption current of the booster circuit, and a soft erase state is established.\nThen, the source voltage automatically rises when the BTBT current is reduced with the progress of erase. By using this method, the consumption current attributed to the BTBT current can reliably be suppressed by the required amount. If the consumption current due to the BTBT current is desired to be suppressed to a half, it is proper to limit the consumption current of the booster circuit to a half.\nHowever, since the source voltage during erase is not specified by this method, it is very difficult to estimate the required erase time.\nNext, a prior art example of the operation of erasing a plurality of blocks will be described with reference to FIGS. 14 and 15. FIG. 14 shows only a part that belongs to a nonvolatile semiconductor storage device capable of executing block erase and is necessary for explaining the erase of the plurality of blocks. A block 1070 and a block 1075, which are the erase blocks, have respective memory arrays 1010 and 1015 and the peripheral circuits thereof. That is, the block 1070 is constructed of a memory array 1010 that becomes an erase unit, a row decoder 1020 for controlling a row select line group 1021 of the memory array 1010, a column decoder 1030 for controlling a column select line group 1031 of the memory array 1010 and a source control circuit 1040 for controlling a source line 1041 of the memory array 1010.\nLikewise, the block 1075 is constructed of a memory array 1015 that becomes an erase unit, a row decoder 1025 for controlling a row select line group 1026 of the memory array 1015, a column decoder 1035 for controlling a column selection group 1036 of the memory array 1015 and a source control circuit 1045 for controlling a source line 1046 of the memory array 1015.\nFurthermore, a debooster circuit 1060 shown in FIG. 14 supplies a negative voltage 1001 to the row decoders 1020 and 1025 provided for a plurality of blocks 1070 and 1075, while a booster circuit 1050 supplies a high voltage 1002 to the source control circuits 1040 and 1045 of the blocks 1070 and 1075. In this case, the debooster circuit 1060 and the booster circuit 1050 have a construction similar to that of the aforementioned booster circuit shown in FIG. 11.\nThe row select line groups 1021 and 1026 shown in FIG. 14 include all the row select lines connected to the memory cells included in the memory arrays 1010 and 1015. The column select line groups 1031 and 1036 include all the column select lines connected to the memory cells included in the memory arrays 1010 and 1015.\nThese memory arrays 1010 and 1015 are the two of the plurality of memory arrays included in this semiconductor storage device. When the storage contents of either one of these memory arrays 1010 and 1015 are erased, the storage contents of a single memory array are erased according to one example by the operation already explained with reference to FIG. 12.\nIn this case, a prior art example of the erase voltage applying method when both the memory arrays 1010 and 1015 are concurrently erased will be described with reference to FIG. 15.\nBetween a time t0 and a time t1, an erase voltage (xe2x88x9210 V, for example) is applied to the row select line group 1021, and an erase voltage (5 V, for example) is applied to the source line 1041. Between the time t1 and a time t2, an erase voltage is applied to the row select line group 1026, and an erase voltage is applied to the source line 1046. As a consequence of this series of operations, the storage contents of the memory cells connected to the memory array 1010 are erased between the time t0 and the time t1, and the storage contents of the memory cells connected to the memory array 1015 are erased between the time t1 and the time t2. Therefore, the storage contents stored in the two memory arrays 1010 and 1015 are entirely erased.\nAlthough the case of the erase of the two blocks 1070 and 1075 has been described with reference to FIGS. 14 and 15, the erase can be executed by a method similar to the aforementioned method when erasing three or more blocks. According to the erase method described with reference to FIG. 15, the time of erase is roughly proportional to the number of blocks erased.\nThen, one example of the method for erasing two blocks at a higher speed will be described with reference to FIG. 16. As shown in FIG. 16, between the time t0 and the time t1, the row line voltage during erase is applied to both of the row select line groups 1021 and 1026, and a source line voltage during erase is applied to both of the source lines 1041 and 1046. By this operation, between the time t0 and the time t1, the storage contents of the memory cells included in both the memory arrays 1010 and 1015 are erased. Therefore, according to the erase method shown in FIG. 16, the erase can be executed at a higher speed than in the case of the aforementioned erase method shown in FIG. 15. Although the case of the erase of the two blocks has been described with reference to FIG. 16, the erase voltage can similarly be simultaneously applied to three or more blocks.\nHowever, the erase method shown in FIG. 16 has the problems as follows. That is, the first problem is that the sum total of the BTBT current flowing from the sources of the memory cells is increased by an increase in the number of memory cells to be simultaneously erased, and when the sum of the currents exceeds the current supply capacity of the booster circuit 1050, the potential of the high voltage 1002 cannot be maintained, hindering the erase operation. The second problem is that the number of blocks that can be simultaneously erased is limited also by a resistance that accompanies the high voltage 1002. The limitations will be described with reference to FIG. 18.\nFIG. 18 shows a model of the source line when four blocks 1170a, 1170b, 1170c and 1170d are provided as a concrete example of the nonvolatile semiconductor storage device that has a plurality of blocks. The blocks 1170a, b, c and d are each provided with a memory array and its peripheral circuits, which are collectively shown in the form of blocks in FIG. 18. Moreover, a circuit for applying a voltage to the row line group and a control circuit are required for the erase operation. FIG. 18 shows elements necessary for explaining the limitations on the number of erase blocks due to the resistance of the source line.\nAs shown in FIG. 18, this nonvolatile semiconductor storage device includes a booster circuit 1150 for generating a source voltage during erase, four blocks 1170a through 1170d that become an erase unit and the resistances 1180a through 1180d of the source lines. This booster circuit 1150 generates a high voltage 1102 and supplies high voltages 1102a through 1102d to the blocks 1170a through 1170d. The resistances 1180a through 1180d of the source lines are caused by the wiring of the node, switching elements and so on and do not normally become zero. The parasitism of the source resistance also changes depending on the arrangement of the blocks, and the case of resistances parasitic in series will be described as the most comprehensible example.\nIn this case, as one example for simplicity, it is assumed that the resistances 1180a through 1180d are all 10 xcexa9, each of the blocks 1170a through 1170d consumes a current of 10 mA at a maximum from the source during erase, the source voltage during erase is permitted to drop to 0.4 V from the voltage generated by the booster circuit 1150, and the current supply capacity of the booster circuit 1150 is sufficient.\nFor example, when the erase of the block 1170a is executed, the currents flowing through the resistances 1180a through 1180d are 10 mA, 0 mA, 0 mA and 0 mA, respectively. A voltage drop due to a resistance becomes the product of the resistance value and the current value, and therefore, the voltage 1102a becomes a voltage dropped by 0.1 V from the voltage 1102. Since this drop voltage is within a permissible range, the erase can be achieved. Likewise, with regard to the blocks 1170b through 1170d, the voltages 1102b, 1102c and 1102d applied to the sources become the voltages dropped by 0.2 V, 0.3 V and 0.4 V, respectively, from the voltage 1102. The voltages are within a permissible range.\nThe case where the four blocks of the blocks 1170a through 1170d are required to be all erased in order to execute the erase of a plurality of blocks (i.e., multi-block erase) at high speed will be described next. If an erase voltage is simultaneously applied to all of the blocks 1170a through 1170d, the currents flowing through the resistances 1180a through 1180d become 40 mA, 30 mA, 20 mA and 10 mA, respectively. Consequently, the source voltages 1102a through 1102d drop from the high voltage 1102 to 0.4 V, 0.7 V, 0.9 V and 1.0 V. The drop of the source voltages of the blocks 1170b through 1170d described above exceeds the permissible range and possibly hinders the erase operation. Therefore, even if the booster circuit 1150 has a sufficient current supply capacity, it is not allowed to simultaneously apply the erase voltage to the four blocks 1170a through 1170d. \nFurthermore, when the erase voltage is simultaneously applied to the two blocks of the aforementioned four blocks, there inevitably occurs a situation that the voltage drop on the source line exceeds the permissible range no matter which combination is selected. Therefore, when erasing these four blocks 1170a through 1170d, there has been repeated four times the erase operation of one block.\nThis source resistance can be reduced by reducing the sheet resistance of the wiring layer of the source line, shortening the source line wiring and increasing the capacity of the switching element when the switching element is employed. As a method for reducing this sheet resistance, the reduction can be achieved by devising the material and thickness of the wiring layer. However, the devising, which requires a change of process, is generally not easy. Moreover, shortening the source wiring has limitations attributed to the layout arrangement of the blocks and so on. Moreover, thickening the source wiring width and increasing the switching element capacity lead to an increase in the chip area, and there is a concern about the sacrifice of the chip size and cost. Therefore, it is unavoidable that a resistance is parasitic on the source line to a certain extent, and this parasitic resistance possibly limits the simultaneous erase of a plurality of blocks.\nAs described above, there have been limitations on the simultaneous erase of a plurality of blocks due to the current supply capacity of the booster circuit that generates the high voltage applied to the source and the resistance parasitic on the source line.\nAs described above, according to the conventional source erase method, the peak value of the consumption current during erase is large, and there is a serious hindrance to the use of the method in portable equipment whose current supply capacity is not sufficiently large and the like.\nMoreover, the substrate erase, which is the erase method that causes no BTBT current, requires substrate potential control, and therefore, a special process and control must be prepared.\nMoreover, according to the method for reducing the peak current by soft erase, it is generally difficult to quantitatively estimate the correlation between the applied source voltage and the peak current value. Using the technology of limiting the consumption current of the booster circuit is also effective for the suppression of the consumption current, but it is difficult to estimate the time necessary for erase.\nMoreover, it has conventionally been difficult to execute batch erase of a plurality of blocks, and when only a single block can be simultaneously erased, the time necessary for erasing a plurality of blocks has been increased roughly in proportion to the number of blocks in which the erase is executed.\nAccordingly, the object of the present invention is to provide a nonvolatile semiconductor storage device capable of quantitatively estimating the required erase time and the peak value of an erase current, suppressing an increase in the erase time and reducing the erase current and a storage contents erase method therefor. The present invention also provides a method capable of efficiently executing the erase of a plurality of blocks.\nIn order to achieve the above object, according to the present invention, there is provided a nonvolatile semiconductor storage device having a memory array constructed of a plurality of memory cells, a row select line control circuit, a column select line control circuit and a source line control circuit for controlling row select lines, column select lines and source lines respectively connected to each of the memory cells and a voltage generation circuit for individually generating voltages to be applied to the row select line and the source line, the memory array being subjected to erase by applying specified voltages to the row select line and the source line, wherein\nthe row select line control circuit comprises a line group independent control means capable of independently controlling a plurality of row select line groups each of which is constructed of at least one row select line, and\nthe voltage generation circuit comprises a consumption current limiting means for limiting a consumption current of the voltage generation circuit.\nIn this invention, the line group independent control means independently controls a plurality of row select line groups, and voltages generated by the voltage generation circuit are applied to the individual row select line groups with time shifts. As a result of this, the peak of consumption current can be suppressed, so that the consumption current can be reduced.\nFurther, in this invention, the consumption current limiting means of the voltage generation circuit limits the consumption current of the voltage generation circuit, allowing voltages to be generated within a range under a specified current value according to the conditions of voltage application from the voltage generation circuit to the individual row select line groups. Thus, a further reduction in consumption current at a shorter scale can be achieved.\nIn one embodiment, there is provided a nonvolatile semiconductor storage device having a memory array constructed of a plurality of memory cells, a row select line control circuit, a column select line control circuit and a source line control circuit for controlling row select lines, column select lines and source lines respectively connected to each of the memory cells and a voltage generation circuit for individually generating voltages to be applied to the row select line and the source line, the memory array being subjected to erase by applying specified voltages to the row select line and the source line, wherein\nthe source line control circuit comprises an independent control means capable of independently controlling a plurality of source lines, and\nthe voltage generation circuit comprises a consumption current limiting means for limiting a consumption current of the voltage generation circuit.\nIn this one embodiment, the independent control means independently controls a plurality of source lines, and voltages generated by the voltage generation circuit are applied to the individual source lines with time shifts. As a result of this, the peak of consumption current can be suppressed, so that the consumption current can be reduced.\nFurther, in this embodiment, the consumption current limiting means of the voltage generation circuit limits the consumption current of the voltage generation circuit, allowing voltages to be generated within a range under a specified current value according to the conditions of voltage application from the voltage generation circuit to the individual source lines. Thus, a further reduction in consumption current at a shorter scale can be achieved.\nIn another embodiment, the consumption current limiting means is constructed of:\nan output current limiting means.\nIn this embodiment, since the consumption current limiting means is implemented by the output current limiting circuit, the voltage generation means outputs voltages within a range under a specified current by virtue of operations of this limiting circuit. As a result of this, output current load is limited, drops due to load currents are reduced, and the power supply capability of the voltage generation means is alleviated, so that the consumption current is limited. This function of reducing the consumption current enables the voltage generation means to be optimum for the source voltage generation circuit in block erase operations in rewritable nonvolatile semiconductor storage devices.\nIn an embodiment, the consumption current limiting means is constructed of:\nan input current limiting circuit.\nIn this embodiment, since the consumption current limiting means is implemented by an input current limiting circuit, the consumption current can be reduced by limiting the current of the power supply current of the voltage generation means.\nIn another embodiment, the consumption current limiting means activates only partially the voltage generation circuit.\nIn this embodiment, since the consumption current limiting means activates only partially the voltage generation circuit, the consumption current of the voltage generation circuit is reduced. This function of reducing the consumption current enables the voltage generation circuit to be optimally usable for the source voltage generation circuit in block erase operations.\nIn an embodiment, there is provided a storage contents erase method for a nonvolatile semiconductor storage device comprising the step of:\nchanging a number of row select line groups to which a specified voltage necessary for erase is simultaneously applied by the line group independent control means according to a specified condition before the erase of all the memory cells in the memory array is completed among a plurality of row select line groups capable of being independently controlled.\nIn this embodiment, the number of row select line groups to which a specified voltage required for erase is simultaneously applied is changed over according to predetermined conditions before the erase of all the memory cells within the memory array is completed. As a result of this, a reduction in peak value of the consumption current can be achieved.\nIn another embodiment, there is provided a storage contents erase method for a nonvolatile semiconductor storage device comprising the step of:\nchanging a number of source lines to which a specified voltage necessary for erase is simultaneously applied by the independent control means according to a specified condition before the erase of all the memory cells in the memory array is completed among a plurality of source lines capable of being independently controlled.\nIn this embodiment, among a plurality of source lines that can be controlled independently by the independent control means, the number of source lines to which a specified voltage required for erase is simultaneously applied is changed over according to predetermined conditions before the erase of all the memory cells within the memory array is completed. As a result of this, a reduction in peak value of the consumption current can be achieved.\nIn an embodiment, the number of the row select line groups or the source lines to which the specified voltage necessary for the erase is simultaneously applied is changed every time a consumption current in the voltage generation circuit for generating the specified voltage becomes not higher than a specified value.\nIn this embodiment, a reduction of peak value of the consumption current can be achieved.\nIn another embodiment, the number of the row select line groups or the source lines to which the specified voltage necessary for the erase is simultaneously applied is changed every predetermined time.\nIn this embodiment, a reduction in peak value of the consumption current can be achieved.\nIn an embodiment, the number of the row select line groups or the source lines to which the specified voltage necessary for the erase is simultaneously applied is changed every time a threshold voltage of an objective memory cell to be erased becomes not higher than a specified value.\nIn this embodiment, the number of row select line groups or source lines to which the erase voltage is applied is changed over stepwise in an erase operation. Thus, the peak consumption current can be reduced efficiently.\nIn another embodiment, overerase verify after an erase operation is executed after completing the erase operation of all objective memory cells to be erased.\nIn this embodiment, after the erase operation of all the memory cells that are targeted for erase is completed, an overerase verify subsequent to this erase operation is executed. Therefore, the memory cells never undergo any erase disturb after the overerase verify. As a consequence, it can be ensured that the threshold voltages of the memory cells fall within a reference threshold voltage range that has been checked at the time of the overerase verify. That is, according to this embodiment, a verify precision equivalent to that of the conventional batch erase method can be obtained.\nIn an embodiment, erase is executed by switchover between an erase method for erasing in a batch a block of a memory cell area selected by the row select line group and the source line by controlling all the row select line groups in an identical operation and controlling all the source lines in an identical operation and the storage contents erase method.\nWith the erase method of this embodiment, it becomes possible to change over between the conventional block erase and the above-described erase method for reducing the consumption current peak, in accordance with the current supply capability of the power supply.\nIn another embodiment, a nonvolatile semiconductor storage device further comprises an erase method switchover means for executing the erase by switchover between an erase method for erasing in a batch a block of a memory cell area selected by the row select line group and the source line by controlling all the row select line groups in an identical operation and controlling all the source lines in an identical operation and the storage contents erase method.\nIn the nonvolatile semiconductor storage device of this embodiment, it becomes possible to change over between the conventional block erase and the above-described erase method for reducing the consumption current peak, in accordance with the current supply capability of the power supply.\nIn an embodiment, there is provided a nonvolatile semiconductor storage device storage contents erase method for executing erase of storage contents of a nonvolatile semiconductor storage device having a plurality of memory blocks including a memory array constructed of a plurality of memory cells, a row select line control circuit, a column select line control circuit and a source line control circuit for respectively controlling row select lines, column select lines and source lines respectively connected to each of the memory cells in each of the memory blocks and a voltage generation circuit for individually generating voltages to be applied to the row select line and the source line, by applying specified voltages to the row select line and the source line, comprising the steps of:\nconcurrently selecting the row select lines connected to the memory block by the row select line control circuit;\nconcurrently selecting the source lines connected to the memory block by the source line control circuit; and\nchanging according to a specified condition a number of the memory blocks to the row select lines of which a specified voltage necessary for erase is simultaneously applied until the erase of all the memory cells included in the memory block in which the erase is executed is completed.\nIn this embodiment, the number of memory blocks to row select lines of which a specified voltage required for erase is simultaneously applied is changed over according to predetermined conditions. Thus, while increases in erase time are suppressed, a reduction in the consumption current can be achieved, so that the erase of a plurality of blocks can be executed efficiently.\nIn an embodiment, the number of the memory blocks to the row select lines or the source lines of which the specified voltage necessary for the erase is simultaneously applied is changed every time a consumption current in the voltage generation circuit for generating the specified voltage becomes not higher than a specified value.\nAccording to this embodiment, while increases in erase time are suppressed, a reduction in the consumption current can be achieved and the peak consumption current can be reduced efficiently, so that the erase of a plurality of blocks can be executed efficiently.\nIn an embodiment, the number of the memory blocks to both the row select line groups and the source lines of which an erase voltage is simultaneously applied is varied with a block in which the connected source line has a comparatively large wiring resistance and with a block in which the connected source line has a comparatively small wiring resistance.\nAccording to this embodiment, while drops of the source voltage in an erase operation are suppressed, a shortening of the erase time can be achieved and a speedup of erase of a plurality of blocks can be achieved.\nIn an embodiment, the blocks to which the erase voltage is simultaneously applied are selected so that a maximum value of a potential drop from an output of the voltage generation circuit for generating a voltage to be applied to the source line inputted to the block in which erase is executed falls within a predetermined permissible range.\nAccording to this embodiment, while drops of the source voltage in an erase operation are suppressed to within a permissible range, a shortening of erase time can be achieved and a speedup of erase of a plurality of blocks can be achieved.\nIn an embodiment, the erase of the memory array in the objective block to be erased is executed by the erase method.\nIn this embodiment, memory arrays within a block targeted for erase can be erased with high efficiency."}
{"text": "HIV-1 is the causative agent of AIDS and presently infects approximately 33 million persons worldwide with approximately 1.9 million infected persons in North America alone. Recent studies have shown that HIV/AIDS has become a global epidemic that is not under control in developing nations. The rapid emergence of drug-resistant strains of HIV throughout the world has placed a priority on innovative approaches for the identification of novel drug targets that may lead to a new class of anti-retroviral therapies.\nThe virus contains a 10-kb single-stranded RNA genome that encodes three major classes of gene products that include: (i) structural proteins (Gag, Pol and Env); (ii) essential trans-acting proteins (Tat, Rev); and (iii) “auxiliary” proteins that are not required for efficient virus replication in permissive cells (Vpr, Vif, Vpu, Nef) [reviewed in (1)]. There has been a heightened interest in Vif as an antiviral target because of the discovery that the primary function of Vif is to overcome the action of a cellular antiviral protein known as APOBEC3G or A3G (2).\nIn 1984, it was determined that HIV was the virus that causes AIDS and researchers declared that a vaccine would be available within two years. Nearly three decades later, there is still no vaccine available and the primary focus remains on developing therapeutics for those already infected. Currently, the primary HIV preventative is the combination treatment known as STRIBILD™. This medication contains 3 components of the highly active anti-retroviral therapy (HAART) regimen (i.e., one integrase (IN) and two reverse transcriptase (RT) inhibitors). However, HIV strains with resistance to some or all of these components had already emerged prior to the availability of STRIBILD™, thus rendering it ineffective against such strains. Furthermore, not only has HIV developed resistance to STRIBILD™ components, but it also has developed resistance to all HAART medications to date, including inhibitors of all HIV enzymatic and viral entry targets. In fact, it is common to see drug resistance even among treatment-naïve individuals worldwide, emphasizing that at least some of the current drugs have limited efficacy in a subset of untreated, infected individuals. The barrier to developing resistance to HIV drugs is low and often a single codon change in the targeted protein is sufficient to cause resistance to more than one inhibitor of the same class (i.e., M46I/L/V in the HIV protease confers resistance to 7 out of 8 inhibitors). The ever-present problem of drug resistance together with the lack of success in developing a vaccine accentuate the need for novel HIV prevention and treatment strategies that are unlikely to develop resistance.\nThe present invention is directed toward overcoming these and other deficiencies in the art."}
{"text": "Food preparation that involves cutting, slicing and chopping is conveniently done over the kitchen sink. Food scraps to be discarded are easily swept into the sink basin with access to the waste disposal unit. Food in stages of preparation is easily washed. For this reason some over-the-sink chopping boards have been suggested. See, for example, U.S. Pat. No. 3,625,162 to Crew, and U.S. Pat. No. 4,765,603 to Huppert as well as U.S. Pat. No. 4,243,184 to Wright. Such devices also add food preparation surface to the existing kitchen counter space which might already be crowded as during preparation of a large meal.\nBroadly stated, the present invention is directed to a cutting board combination for use with a double basin sink of the type having a sink divider between the basins. The combination cutting board includes a first, smaller cutting board adapted to be removably mounted on the divider separating the sink basins. The first cutting board is comprised of a generally rectangular hard slab-like member having a top, flat, relatively indestructible smooth work surface. A pair of spaced apart vertically depending legs are rigidly secured to the bottom surface of the first cutting board. The legs are adapted to straddle the sink divider with the cutting board resting upon it whereby the first cutting board is usable as a stand-alone cutting board with lateral edges that overhand the sink basins.\nThe cutting board combination includes a second, larger cutting board that can be connected to the first cutting board to increase the work surface area. One lateral edge of the second cutting board abuts a lateral edge of the first cutting board. The second cutting board is long enough to span the width of the sink basin so that the opposite lateral edge overlaps the sink rim. The second cutting board is likewise comprised of a flat, rectangular slap-like member with a flat top surface that is preferably flush with the top surface of the first cutting board when the two are connected. The first cutting board has connecting means on both lateral edges so that the second cutting board can be connected to either side to span either sink basin."}
{"text": "In general, this invention relates to an ablatable laminar imaging medium having an optimized exposure threshold, and in particular, one wherein said optimized threshold is resultant of the employment therein of an ablation enhancement layer with predefined absorption and reflection values.\nLiquid crystal displays comprise a liquid crystal material sandwiched between two substantially transparent electrode assemblies. Touch screen displays of either the resistive or capacitive types comprise a display screen (for example, a cathode ray tube) having superposed thereover two substantially transparent electrode assemblies. In both types of displays, each of these electrode assemblies typically comprises a substrate on which is deposited a conductive layer thin enough to be substantially transparent.\nThe term xe2x80x9csubstantially transparentxe2x80x9d is used herein to mean that the electrodes transmit sufficient visible light so that the two superposed electrodes will not substantially obscure, nor substantially distort the color of, a liquid crystal display or touch screen display incorporating the two electrodes. Typically, commercial specifications require that the two superposed electrodes have a transmittance of at least 80% at 550 nm.\nIn liquid crystal displays the substrate is usually glass, whereas touch screen displays usually employ a synthetic resin (plastic) substrate for at least one electrode. The conductor is often formed from indium tin oxide or a similar metal oxide. The conductor is typically formed by depositing the oxide by sputtering or chemical vapor deposition at a high temperature, and then annealing, also at a high temperature. On glass substrates temperatures in excess of 300xc2x0 C. may be used to deposit and anneal the conductor; on plastic substrates, lower temperatures must be used, with resultant higher electrical resistance in the conductor.\nAlternatively, both liquid crystal displays and touch screen displays may make use of thin film electrodes comprising a metallic conductive layer sandwiched between two layers having high refractive indices; these two layers usually being formed from metal oxides. The metallic conductive layer is patterned so as to divide it into a plurality of electrodes, and conductors are attached to each of these electrodes to enable formation of the desired patterns in the liquid crystal material.\nConventional prior art (i.e., photolithographic) processes for forming electrodes generally involve the deposition of layers of photoreactive and non-photoreactive thin films (typically, metal oxides) onto a substrate, the selective exposure of said layers through a mask (or like phototool) corresponding to the desired electrode pattern, and removal of either the exposed or unexposed portionsxe2x80x94depending on the nature of the photoreaction involvedxe2x80x94of the thin film layers.\nPrior art processes for forming electrodes also often require the use of elevated temperatures of 200xc2x0 C. or more, which in practice requires the use of glass substrates or expensive high temperature plastics (polymers are known which have glass transition temperatures above 225xc2x0 C. and can thus withstand processing at such temperatures). There are many applications for liquid crystal displays (for example, in cellular telephones and other mobile electronic devices) where it would be advantageous to use less expensive plastic substrates having lower glass transition temperatures if thin film electrodes could be formed on such substrates.\nWhile applicability of the above processes still remain practical and commercially viable, interest in the formation of electrode patterns by direct laser imaging (particularly, laser ablation processes) continues to develop at a relatively high ratexe2x80x94in part, because of laser imaging\"\"s potentially higher accuracy, its faster production, and its independence from certain costly pre-imaging processes (e.g., mask preparation). Potential also exists in the use of lighter and cheaper plastic substrates.\nThe formation of LCD electrode patterns by laser ablation essentially involved the direct pattern-wise removal of portions layers of electrically conductive and insulating materials by exposing said portions to laser light of an intensity and quality (e.g., wavelength) sufficient to completely or partly decompose said material. The reaction can be fairly characterized as xe2x80x9cexplosivexe2x80x9d (i.e., on a microscopic level), producing vapor-like or gas-like streams consisting of fragments of the removed material. An example of such process is disclosed in U.S. patent application Ser. No. 09/009,391, filed by Hyung-Chul Choi, Yi-Zhi Chu, Linda S. Heath, and William K. Smyth on Jan. 20, 1998.\nOverall, most laser ablation electrode-forming processes produce good results. Regardless, areas for further improvement exist.\nFor example, it has been observed that the process of ablating a metal and partial conductor sometimes results in the formation of ridges (like undesirable physical anomalies) as a result of the volcanic action of ablated metal from beneath the partial conductor. These ridges reduce boundary definition of ablated areas, and as such, can potentially compromise the pixel resolution in a finished flat panel display. While conventional post-patterning cleaning processes can be employed to address this problem, the debris on etched edges is difficult to remove, and requires vigorous steps that can damage the thin film structure.\nAlso, in the conduct of conventional laser ablation processes, high threshold exposure energies are typically involved. Aside from the apparently high energy costs, use of intense ablation exposure produces correspondingly intense temperature elevations in and around the area of ablation. This limits selection of material used for the ablation medium to those having high thermal resistance. And, in the selection of substrates, this greatly constrains the use of plastic materials, which is often preferred over the more common glass substrates when lower costs, greater flexibility, and lighter weight are desired.\nFinally, laser ablation produces residue, which needs to be washed off. Being essentially composed of the same material that ultimately becomes the finished electrode, such residue, if left in the valleys between electrodes, can cause an electrical short. The presence of such a short is of course undesirable, since it in effect turns the two adjacent electrodes into a single electrode and thereby adversely affects the quality of a liquid crystal or touch screen in which the electrode assembly is used.\nIn response to the above-identified issues, an ablatable laminar imaging medium is described herein, characterized by its incorporation of an ablation enhancement layer, the ablation enhancement layer being composed, configured, and located to effectively reduce the threshold energy requirement for effecting laser ablation. The use of the ablation enhancement layer dramatically improves laser ablation efficiency (i.e., faster scan speed, lower exposure energy threshold, and quality), and promotes xe2x80x9ccleanerxe2x80x9d ablation of the metallic layer(s) in the area of exposure (e.g., well ablated areas with few xe2x80x9cridgesxe2x80x9d, unablated residue, and/or other like ill-defined edges, surfaces, and boundaries).\nIn particular, the ablatable laminar imaging mediumxe2x80x94which is particularly useful in the manufacture of a substantially transparent electrode assemblyxe2x80x94can be defined as comprising: (a) a substrate; (b) a high-index metal oxide layer; (c) an ablatable metallic conductive layer; (d) a high-index conductive metal oxide layer; and (e) an ablation enhancement layer having an IR-absorption greater than the IR-absorption of said high-index conductive metal oxide layer and an IR-reflectivity less than the IR-reflectivity of said high-index conductive metal oxide layer.\nThe present invention also provides a new process for dealing with (cf., xe2x80x9ccleaningxe2x80x9d) remnant and/or residual ablation enhancement material left in the course of ablation, which for many applications and products is unwanted. The characterizing steps of the inventive method involve the execution of a xe2x80x9ccleaningxe2x80x9d operation in close connection with a xe2x80x9cpatterningxe2x80x9d operation, preferably using the same apparatus (e.g., the same laser) for both. The new process can be accomplished by either a sequential methodology (i.e., ablation patterning, followed by ablation cleaning) or a continuous methodology (i.e., ablation patterning conducted contemporaneously with ablation cleaning ).\nIn light of the above, it is a principal object of the present invention to provide an ablatable laminar imaging, useful in the manufacture of an electrode assembly for a flat panel display, having a low exposure threshold requirement.\nIt is another object of the present invention to provide an ablatable laminar imaging medium having incorporated therein an ablation enhancement layer, the IR-reflectivity and IR-absorptivity of the ablation enhancement layer specifically designed for favorable effect on the medium\"\"s exposure threshold.\nIt is another object of the present invention to provide photo-ablation means for making an electrode assembly that requires less exposure energy, thus generating less heat, and thus allowing for use in manufacture of relatively inexpensive, heat-sensitive stacks and light-weight substrates (such as plastic substrates) that typically have less thermal resistance than glass.\nIt is another object of the present invention to provide a method for patterning an ablatable laminar imaging medium, wherein the patterning operation is executed in close association with a cleaning operation, ie., the cleaning operation being accomplished by laser ablation after the patterning operation, preferably by the same apparatus used therefor.\nIt is another object of the present invention to provide a method for patterning an ablatable laminar imaging medium, wherein the patterning operation is executed in close association with a cleaning operation, ie., the cleaning operation being accomplished by laser ablation contemporaneously with the patterning operation, preferably by the same apparatus used therefor.\nOther objects of the present invention will become apparent from the detailed description below taken in conjunction with the accompanying drawings."}
{"text": "One of the problems which is encountered in target shooting with a pistol is the human reaction which causes a shooter to flinch or slightly move his hand, at or about the time the trigger movement releases the hammer. I have observed that much greater accuracy can be achieved if the hammer travels through a shorter distance and therefore for a shorter period of time, and the main object of this invention is to provide improvements which will achieve this result."}
{"text": "Embedded controllers (also called embedded devices, embedded controllers, embedded systems or, simply, microprocessors) generally are devices whose primary purpose is to perform an independent action such as encoding video, controlling robots, routing internet traffic, and so on. Such devices run autonomously without user intervention, however, there needs to be a way to configure/monitor/troubleshoot these devices.\n“Embedded systems” are a type of systems (often including a plurality of embedded devices in one product that is offered for sale) that have applications that that are growing rapidly and expanding into different markets. One common typical characteristic of devices in this group is “fixed functionality”. These devices start with a predefined set of application tasks and this set cannot be modified. These devices do not allow users to upload and run third-party programs. This assumption allows the manufacturer to reduce software complexity and to create very small, inexpensive and reliable systems.\nThe current generation of conventional embedded networking devices use standard off-the-shelf components such as operating systems (OS), file systems, web servers and SNMP agents (Simple Network Management Protocol (SNMP) is an internet-standard protocol for managing devices on IP networks) in order to provide the networking and application tasks. Due to various reasons, these systems typically involve expensive and complex hardware designs and inefficient software designs. Due to such inefficiencies, these systems contain huge amounts of code and require large amounts of RAM and high CPU speeds to support any reasonably robust networking infrastructure. This makes the devices expensive and unsuitable for certain applications. The large amount of code also results in a very slow startup time which blocks the adaptation of networking into certain devices such as TV sets, DVD players and other kinds of devices which need to be turned off and on quite often and require boot-up times on the order of one second. A number of U.S. patents attempt to facilitate the development of embedded controllers (also called embedded devices, embedded controllers, or, simply, microprocessors) in various ways, including the following U.S. Patents, each of which is incorporated herein by reference.\nU.S. Pat. No. 8,209,694 to Glistvain issued on Jun. 26, 2012 with the title “Single-stack real-time operating system for embedded systems”, and is incorporated herein by reference. This patent describes real time operating system (RTOS) for embedded controllers having limited memory includes a continuations library, a wide range of macros that hide continuation point management, nested blocking functions, and a communications stack. The RTOS executes at least a first and second task and uses a plurality of task priorities. The tasks share only a single stack. The task scheduler switches control to the highest-priority task. The continuations library provides macros to automatically manage the continuation points. The yield function sets a first continuation point in the first task and yields control to the task scheduler, whereupon the task scheduler switches to the second task and wherein at a later time the task scheduler switches control back to the first task at the first continuation point. The nested blocking function invokes other blocking functions from within its body and yields control to the task scheduler.\nU.S. Pat. No. 7,383,368 to Schopp issued on Jun. 3, 2008 with the title “Method and system for autonomically adaptive mutexes by considering acquisition cost value”, and is incorporated herein by reference. This patent describes a method for managing a mutex in a data processing system. For each mutex, an average acquisition cost is maintained that indicates an average consumption of computational resources that has been incurred by threads attempting to acquire the mutex.\nU.S. Pat. No. 6,282,454 to Papadopoulos, et al. issued on Aug. 28, 2001 with the title “Web interface to a programmable controller”, and is incorporated herein by reference. This patent describes an interface to a network of at least one programmable logic control system running an application program for controlling output devices in response to status of input devices. The web interface runs web pages from an Ethernet board coupled directly to the PLC back plane and includes an HTTP protocol interpreter, a PLC back plane driver, a TCP/IP stack, and an Ethernet board kernel. The Web interface provides access to the PLC back plane by a user at a remote location through the Internet. The interface translates the industry standard Ethernet, TCP/IP and HTTP protocols used on the Internet into data recognizable to the PLC. Using this interface, the user can retrieve all pertinent data regarding the operation of the programmable logic controller system. This is different than the present invention because U.S. Pat. No. 6,282,454 describes a web interface module and PLC that are separated via a backplane interface, while, in contrast, in the present invention the web interface is implemented inside the PLC (or industrial-automation controller or other microcontroller), not as one or more separate parts.\nU.S. Pat. No. 7,058,693 to Baker, Jr. issued on Jun. 6, 2006 with the title “System for programming a programmable logic controller using a web browser”, and is incorporated herein by reference. U.S. Pat. No. 7,058,693 describes a control system that includes an internet web interface to a network of at least one programmable logic control system (PLC) running an application program for controlling output devices in response to status of input devices. The web interface runs web pages from an Ethernet® board coupled directly to the PLC back plane and includes an HTTP protocol interpreter, a PLC back-plane driver, a TCP/IP stack, and an Ethernet board kernel. The web interface provides access to the PLC back plane by a user at a remote location through the internet. The interface translates the industry standard Ethernet, TCP/IP and HTTP protocols used on the Internet into data recognizable to the PLC. Residing in the PLC is a programming package accessible to a user through this interface which will allow the user to edit programs controlling the operation of the programmable logic controller system. Again, this is different than the present invention because U.S. Pat. No. 7,058,693 describes a web interface module and PLC that are separated via a backplane interface, while, in contrast, in the present invention the web interface is implemented inside the PLC (or industrial-automation controller or other microcontroller), not as one or more separate parts.\nPrior art systems include:                (1) statPLC web, described as browser based programming, wherein one points their favorite web browser to the statPLC web server and starts editing. (See www.autstat.com/index.php/statplc-web.html accessed Aug. 19, 2015, which indicates a postal address at Baier automation & statistics, Ufergasse 68, 3500 Krems, Austria.) There is no precise information about how things are hosted, but based on the available information it is likely that it works in the following way: the user runs a browser application on their local client computer, which causes an application to run on a web-based server computer to edit and compile programs to obtain executable code, that is then loaded from the web-based server into the embedded device local to the user. The fact that statPLC web supports C development likely indicates that they have some sort of C compiler, and it is unlikely that the C compiler would run on Cortex M3 CPU that they claim to be supporting.        (2) Tri-Logic (a Java applet-hosted inside the device because of Java limitation), which supports simulation (See www.tri-plc.com/trilogi3.htm, accessed Aug. 19, 2015, which indicates a postal address at 1685 H ST #198, Blaine Wash. 98230, United States.).        (3) ICPDAS WISE (hosted inside the device). (See wise.icpdas.com/Introduction.html, which indicates its Headquarters at No. 111, Guangfu N. Rd., Hukou Township, Hsinchu County 30351, Taiwan, R.O.C.)        (4) MBLOGIC ladder editor on Web (hosted on the same device). (See mblogic.sourceforge.net/techdemo/ladtest/ladderhelp.html)        (5) Ladder with compiler on the client-side generating code for the target device (Arduino/AVR). (See github.com/tadpol/Avrian-Jump.) This can be combined with monitoring Ladder in real-time using Jayscript (which also provides some simulation). (See mblogic.sourceforge.net/techdemo/techdemo.html.)        (6) Managing local Storage to avoid losing work with programming environment (even managing multiple sessions within the same browser). (See mblogic.sourceforge.net/techdemo/techdemo.html.)        (7) Remote access to devices (See research.microsoft.com/pubs/161386/gibraltar.pdf and www2014.kr/wp-content/uploads/2014/05/companion_p135.pdf, accessed Aug. 19, 2015)        (8) Beck PLC (Web PLC) Hosting is on the device itself, wherein the device allows one to upload information to a central Beck server for IOT functionality. (See www.beck-ipc.com/files/manual/com.tom_WEB-PLC_GettingStarted_V16.pdf)        \nAccording to Wikipedia (en.wikipedia.org/wiki/Ladder_logic, accessed Aug. 19, 2015), “ladder logic” was originally a written method to document the design and construction of relay racks as used in manufacturing and process control. Each device in the relay rack would be represented by a symbol on the ladder diagram with connections between those devices shown. In addition, other items external to the relay rack such as pumps, heaters, and so forth would also be shown on the ladder diagram.\nLadder logic has evolved into a programming language that represents a program by a graphical diagram based on the circuit diagrams of relay logic hardware. Ladder logic is used to develop software for programmable logic controllers (PLCs) used in industrial control applications. The name is based on the observation that programs in this language resemble ladders, with two vertical rails and a series of horizontal rungs between them. While ladder diagrams were once the only available notation for recording programmable controller programs, today other forms are standardized in IEC 61131-3.\nAccordingly, there is a need for improved methods and apparatus for developing software for PLCs."}
{"text": "In the early development of data display systems, it was customary to employ a cathode ray tube wherein a layer of phosphor material was made to luminesce by means of an electron beam that scanned across the layer of phosphor material. Although good quality images can be created in this manner, the physical size of the images that can be created by a cathode ray tube is severly limited by various factors, such as, for example, power required and distortion of the electron beam path. In order to provide images of greater size, numerous schemes have been proposed for optically enlarging an image created by cathode ray tube. Although an image of greater proportions can be obtained in this manner, the amount of light available from the phosphors present on the face of the tube is very limited. As a result, the quality of the enlarged image and particularly the brightness thereof has heretofore been very poor.\nCathode ray tubes have also been utilized in systems which employ a projection system having a light source which is independent of the cathode ray tube light emission. For example, in U.S. Pat. Nos. 3,667,830; 3,701,586; 3,609,222; and 3,746,785, large scale displays are provided by utilizing a display structure having a deformable, light reflective metalic film supported by a support grid which is situated within a cathode ray tube. When an electron beam scans across the display structure, charge accumulates on areas of the display structure in accordance with the information content of the electron beam. This charge accumulation causes small deformations or dimples to form in the metal film at the areas of the charge accumulation. When light from a projection system is directed upon the metal film, only light which strikes the deformed areas reaches a display screen. Thus, a light image is formed on the display screen corresponding to the dimpled image formed in the metal film by the electron beam. A flood gun must be provided within the cathode ray tube to neutralize the charge areas to thereby allow the deformations in the metal film to relax to the non-deformed or normal position.\nThe size limitations of display systems using cathode ray tubes has led to the use of matrix addressed displays when large displays are required. In a matrix display, pairs of conductors in a two-dimensional array of such conductors are utilized to address each elemental area of the display to thereby initiate the emission of light at a selected elemental area when the pair of conductors associated with that selected elemental area are properly biased. in such a display system, as described in U.S. Pat. No. 3,091,876, when a selected elemental area is properly biased, a pressure valve is opened which forces a portion of a flexible membrane out past the end of a tubular support member. In this outer position, the reflective surface of the membrane reflects incident light to provide a visible \"bright spot\" in the surface of the display to thereby provide a visible display. U.S. Pat. No. 3,091,876 also teaches using an electroluminescent panel in conjunction with the flexible membrane whereby when an elemental portion of the membrane is addressed, the membrane is forced by a pressure system into contact with an electrode of the electroluminescent panel whereby a voltage is applied across an elemental area of the electroluminescent panel to thereby cause it to emit visible light.\nAs noted, the display systems described which utilize cathode ray tubes suffer from size limitations and the expensive cathode ray tube component. Also, these systems do not use ambient light as the projection light source. The matrix display systems that utilize a pressure source also suffer from the requirement of expensive mechanical components and also, especially when using an electroluminescent panel, from the lack of a threshold behavior since the elemental areas adjacent a selected elemental area receive half the voltage applied across the selected elemental area and that voltage may be sufficient to initiate undesirable glow discharges and hence undesirable light output at areas adjacent the selected elemental area. Also, many of the display systems described do not have machine readable capabilities."}
{"text": "Locking lid and container configuration have been the subject of many efforts. In addition to the applicant's earlier efforts, many of which have been patented, others such as U.S. Pat. Nos. 3,902,620, 6,176,381, 5,544,768 and others show various locking configurations in which the locking mechanism is secured to the container and engages teeth on the lid. Similar techniques have been employed for providing a tamper indicator and is shown in U.S. Pat. No. 6,926,165. All these prior art designs are various ways of achieving the various objectives, however, the applicant believes that there is still room for improvement in the art of containers with locking lids."}
{"text": "Prior art devices for spraying fragrances, deodorising agents and sanitising fluids into a room generally consist of a device containing a removable source of fluid. With such an arrangement once the source of fluid has been completely exhausted, the source can be replaced rather than replacing the entire device. Typically such sources come in many forms including containers, bottles, canisters and cartridges (generically, all such containers, bottles, cans and cartridges hereinafter will be referred to as “refills”). Such refills can be pump sprays or aerosols, including metered and non-metered versions thereof.\nKnown prior art devices are typically provided with a recess to permit wall mounting and with a flat base to permit the device to be free standing on a flat surface. These devices typically comprise a two-part housing of a main body and a cover, the main body holds an actuation mechanism typically consisting of an electrically powered motor, gears, actuation arm, batteries and the like. The main body is further provided with a refill receiving means which typically is formed of a support plate at or adjacent the base to support the base of a refill. The main body is further provided with a support channel adjacent the actuation mechanism, the support channel being present to laterally support a spray head of the refill during actuation, ie—during depression of the spray head by the actuation arm of the actuation mechanism during operation to spray a quantity of fluid from the refill out into the surrounding environment through the spray head. The cover is provided with a spray opening configured to be adjacent the support channel such that an exit nozzle of the spray head will face the opening to permit, in use, sprayed fluid to exit the device. The main body and cover are connected via a hinge adjacent the base of the device to form a closure via a clamshell arrangement to seal the refill inside the device. There are also examples of prior art where the hinge is provided along the side of the main body to produce the clamshell arrangement.\nDespite the prior art devices being commonplace anecdotal data suggests that upwards of 25% of refills are purchased without the purchaser knowing that such refills are intended for use in an automatic spraying device and/or without the purchaser knowing which automatic spraying device the purchased refill is intended to be used with.\nIt is an object of the present invention to remove, reduce or ameliorate the above-mentioned drawbacks as well as address other concerns with the prior art devices."}
{"text": "Toy track systems are typically configured in form of open or uncovered tracks, so that a player may feel the sense of speed by directly watching the track vehicles moving therein. However, according to the above toy track system, a track vehicle moving therein may fall out of the track when passing certain positions at a high speed, particularly at curved or slanted track portions. The track vehicle falling out at a high speed may be easily damaged or even poses a danger to players or bystanders.\nAlthough a toy track system in form of connected pipes may be able to substantially address the above issue, such track system in the form of pipes substantially encloses the track vehicle moving therein. Thus, the player or the bystander cannot directly feel the speed and thereby greatly hindering enjoyment.\nA typical track system only allows the track vehicle to run on a preset track at all times, and such track vehicle would usually be operated in a passive manner. For example, typical track vehicle has only a manual power switch, hence the track vehicle only goes forward along the preset track. In addition, a typical track system would only allow the player to assemble and erect the track on a flat surface, such as the floor or a table. These types of track systems have very limited variations and thus cannot satisfy some players or bystanders.\nThe present invention aims to eliminate or at least alleviate such deficiencies by providing an innovative and improved toy track system and track vehicles."}
{"text": "1. Technical Field\nExample embodiments of the present invention relate in general to technology for a packet-based mobile communication system, and more particularly, to methods of managing a terminal performed in a base station and the terminal to support mobility management and low-power operation of the terminal in a heterogeneous network (HetNet) environment.\n2. Related Art\nIn a packet-based cellular mobile communication system, mobility management is intended to maintain continuity of a connection of a radio bearer (RB), which is a logical channel established between a base station and a terminal.\nA general mobility management method uses a backward handover procedure in which a terminal establishes a connection for a control channel through preliminary information exchange between a source base station from which the terminal has originally received service and a target base station that will provide the service through a new connection according to movement of the terminal, accesses the target base station to set a connection for a data channel, and then closes a connection with the source base station.\nIn particular, in a HetNet environment in which various types of base stations, cells or transmission points such as macro base stations, micro or pico base stations, home base stations, closed subscriber group (CSG) cells or remote radio heads (RRHs) coexist, ping-pong handover frequently occurs between adjacent base stations/cells, and load of a system increases. Also, the quality of a radio channel deteriorates due to interference at a cell boundary, and radio link failure (RLF) occurs due to the deterioration, so that the overall performance of a mobile communication system is hindered.\nTo improve performance of a mobile communication system, there is a necessity for a method capable of reducing signaling overhead caused by message transmission or state switching (e.g., switches between a connected state and an idle state) of terminals, which are smart phones intermittently generating a small amount of data such as instant message (IM) traffic, machine type communication (MTC) terminals generating traffic such as smart metering data, or terminals generating a small amount of data at long intervals, and minimizing power consumption of the terminals.\nTerminals supporting a device-to-device (D2D) communication function are required to perform a measurement for checking whether or not there are adjacent terminals and whether or not it is possible to communicate, and thus it is necessary to improve a measurement procedure such as measurement setup, measurement execution, and measurement reporting for D2D communication, and a method and a control procedure for measurement triggering.\nMeanwhile, to improve performance of a terminal located at a cell boundary in a packet-based cellular mobile communication system, a coordinated multipoint transmission and reception (CoMP) function in which a plurality of base stations, cells or transmission nodes located at geographically the same point or different points cooperate to provide service is taken into consideration. CoMP may be classified as a CoMP function provided by one or more micro base stations in a macro base station having a relatively large arbitrary service area, a CoMP function provided by two or ore macro base stations, a CoMP function provided by two or more macro base stations and micro base stations belonging to the macro base stations, and so on.\nIn addition, CoMP functions may be classified into a joint processing (JP) scheme in which a plurality of transmission nodes transmit the same packet information all together, and a coordinated scheduling/beamforming (CS/CB) scheme in which a plurality of transmission nodes cooperate to support minimization of interference at one transmission node.\nAccording to the JP scheme, in general, a plurality of transmission nodes transmit the same information using the same radio resources (which means that radio resources consisting of the same frequency band and transmission time are allocated to a CoMP-target terminal) and the same modulation and coding scheme (MCS).\nThe CS/CB scheme is a control method in which a plurality of transmission nodes share radio resources and information for MCS to transmit data to an arbitrary terminal using optimal radio resources and MCS, and exchange related information (e.g, the magnitude of an interference signal, the magnitudes of a signal from a serving cell and a signal of an adjacent cell, optimal transmission and coding information such as a pre-coding matrix indicator (PMI) in consideration of the adjacent cell) so that provision of service and allocation of radio resources can be made through interference control.\nThe JP scheme may be classified into a joint transmission (JT) scheme in view of a downlink from a base station to a terminal, and a joint reception (JR) scheme in view of an uplink from a terminal to a base station. Also, the JP scheme may include a dynamic cell selection (DCS) method or a dynamic point selection (DPS) method of dynamically selecting transmission nodes participating in CoMP.\nThe DCS/DPS method is intended to select an optimal point (or cell) at an arbitrary transmission time from among a plurality of cells or points set to participate in CoMP in consideration of radio channel quality, the load statuses of base stations, transmission/reception power and an interference state between a terminal and the base stations, etc., thereby improving performance.\nIn order to improve performance of a mobile communication system in a HetNet environment supporting a CoMP function in which a plurality of base stations cooperate to provide service as described above and a carrier aggregation (CA) function, it is necessary to improve a measurement operation required for connection control between one or more base stations and a terminal, and a discontinuous reception (DRX) operation control procedure. In addition, to improve the performance of the system, it is necessary to improve a mobility control procedure of a terminal, a measurement method for an MTC terminal, a terminal generating traffic having a variety of profiles, and a terminal supporting the D2D communication function, and a control procedure for minimizing power consumption of a terminal."}
{"text": "There is disclosed in DE-B-38 03 212 (U.S. Pat. No. 4,881,523) an endoscope for nasal surgery having an outer endoscope shaft comprising a suction and flushing shaft and a handle connected to the outer shaft, a switching valve being provided in the handle for connecting the suction and flushing shaft to, and disconnecting it from, a source of negative pressure and a flushing fluid source. An optical system extends centrally through the endoscope shaft, so that a further shaft for receiving a surgical instrument cannot be provided."}
{"text": "Gas turbine engines are widely used to power aircraft throughout the world. The engine provides thrust which powers the aircraft by burning a mixture of fuel and air in one or more combustors. A fuel nozzle sprays such mixture into each combustor in a form suitable for rapid mixing and efficient combustion.\nThe most common types of fuel nozzles use a pressure atomizing principle to provide a uniform distribution of fine fuel particles, or droplets, throughout the range of fuel flow conditions encountered during engine operation. In order to be commercially useful, fuel nozzles must be able to (a) efficiently atomize fuel at low air flow rates, (b) uniformly atomize fuel at high power regimes, and (c) provide predictable and controllable fuel spray characteristics over a range of engine operating conditions. Those skilled in the art recognize that other characteristics of fuel nozzles are also desired in addition to those enumerated above.\nWhile progress has been made in designing fuel nozzles for gas turbine engine use, further improvements are required. The present invention provide such improvements."}
{"text": "The desire to expand the frontiers of computer science has prompted consideration of the factors that contribute to the limitations of current computers. Silicon is at the heart of today's computer. The advances in computing power and speed over the years have largely been a consequence of better understanding the fundamental properties of silicon and harnessing those properties for practical effect. Initial progress was predicated on building basic electronic components such as transistors and diodes out of silicon and later progress followed from the development of integrated circuits. More recent advances represent a continuation of these trends and currently emphasize miniaturization and the integration of an ever larger number of microelectronic devices on a single chip. Smaller devices lead to higher memory storage densities, more highly integrated circuits and reduced interaction times between devices on the same chip.\nSince future improvements in computing power and functionality are currently predicated on further improvements in silicon technology, there has been much recent discussion about the prognosis for continued miniaturization of silicon-based electronic devices. A growing consensus is emerging that believes that the computer industry is rapidly approaching the performance limits of silicon. The feature size in today's manufacturing technologies is approximately 0.10 micron and it is expected that this can be reduced to about 0.02 micron in the future. Further decreases in feature size, however, are deemed problematic because sizes below about 0.02 micron lead to a change in the fundamental behavior of silicon. More specifically, as the dimensions of silicon devices decrease to tens of nanometers and below, silicon enters the quantum regime of behavior and no longer functions according to the classical physics that governs macroscopic objects. In the quantum regime, phenomena such as tunneling lead to delocalization of electrons across many devices. Consequences of tunneling include leakage current as electrons escape from one device to neighboring devices and a loss of independence of devices as the state of one device influences the state of neighboring devices. In addition to fundamental changes in the behavior of silicon, further decreases in the dimensions of silicon devices also pose formidable technological challenges. New and costly innovations in fabrication methods such as photolithography will be needed to achieve smaller feature sizes.\nOne strategy for advancing the capabilities of computers is to identify materials other than silicon that can be used as the active medium in data processing and/or storage applications. Such alternative computing media could be used independent of or in combination with silicon to form the basis of a new computing industry that seeks to offer better performance and more convenient manufacturing than is possible with silicon.\nChalcogenide materials are an emerging class of alternative materials for the storage and processing of information. Chalcogenide materials have been previously utilized in optical and electrical memory and switching applications and some representative compositions and properties have been discussed in U.S. Pat. Nos. 5,543,737; 5,694,146; 5,757,446; 5,166,758; 5,296,716; 5,534,711; 5,536,947; 5,596,522; and 6.087,674; the disclosures of which are hereby incorporated by reference herein, as well as in several journal articles including “Reversible Electrical Switching Phenomena in Disordered Structures”, Physical Review Letters, vol. 21, p. 1450-1453 (1969) by S. R. Ovshinsky; “Amorphous Semiconductors for Switching, Memory, and Imaging Applications”, IEEE Transactions on Electron Devices, vol. ED-20, p. 91-105 (1973) by S. R. Ovshinsky and H. Fritzsche; the disclosures of which are hereby incorporated by reference herein.\nChalcogenide phase-change materials form the basis of OUM (Ovonic Universal Memory) technology. OUM is a non-volatile form of memory that is viewed in the near term as a viable alternative to flash memory and DRAM and in the longer term as a viable alternative to SRAM. The functional characteristic of chalcogenide phase-change materials that underlies memory operation is the ability of phase-change materials to undergo a reversible transformation between two or more structural states. The chalcogenide phase-change materials have structural states that include a crystalline state, one or more partially crystalline states and an amorphous state. The crystalline state may be a single crystalline state or a polycrystalline state. A partially crystalline state is a structural state of a phase-change material that includes an amorphous portion and a crystalline portion. Chalcogenide phase-change materials include a plurality of partially crystalline states that differ in the relative proportion of the amorphous and crystalline portions included within a volume of the material. The various structural states of a phase-change material may be distinguished on the basis of electrical resistance. Memory functionality can be achieved by associating different memory states with different structural states and using electrical resistance as the means to read the memory device and discriminate among the different memory states. In a binary memory device having memory states “0” and “1”, for example, state “0” may be associated with a substantially crystalline state and the state “1” may be associated with a substantially amorphous state. Since the resistance of a substantially crystalline state is at least an order of magnitude lower than the resistance of a substantially amorphous state, the two states are readily distinguished through a simple resistance measurement. The operation of storing information (writing or programming) occurs by providing energy (most commonly in the form of electrical current pulses) to the phase-change material to induce the structural transformations needed to establish the desired proportions of crystalline and amorphous phase domains within a volume of the phase-change material. Controlled applications of energy can be used to reversibly and continuously vary the relative proportions of crystalline and amorphous phase domains to establish the structural state corresponding to the information that the programmer wishes to store. Once established, a memory state is stable until further energy having a magnitude sufficient to reprogram the material is applied. The current used to determine the resistance of the device (and thus to read the device) is too low to alter the structural state of the phase-change material.\nChalcogenide switching materials form the basis of the Ovonic Threshold Switch (OTS) technology. Chalcogenide switching materials are substantially amorphous materials that exhibit little or no tendency to undergo a structural transformation to a crystalline or partially crystalline state, but which instead undergo rapid switching from a resistive state to a conductive state upon application of a threshold voltage, Vth. According to a leading model of the switching event, application of the threshold voltage causes the formation of a conductive channel or filament within the chalcogenide material. At the threshold voltage, the electric field experienced by the material is sufficiently, high to induce a breakdown or avalanche-like effect whereby electrons are removed from atoms to form a highly conductive, plasma-like filament of charge carriers. Rather than being bound to atoms, some electrons become unbound and highly mobile. As a result, a conductive channel or filament forms. The conductive filament constitutes a conductive volume within the otherwise resistive chalcogenide material. The conductive filament extends through the chalcogenide material between the device terminals and provides a low resistance pathway for electrical current. Creation of a conductive state upon switching enables the device to support high currents.\nIn order to advance a new chalcogenide-based computing paradigm, it is necessary to develop devices and circuits for performing data storage and processing operations. Chalcogenide OUM technology provides a versatile and robust memory platform for storing data. Representative examples of the application of chalcogenide phase change materials to data processing include mathematical operations (U.S. Pat. No. 6,671,710 (“Methods of Computing with Digital Multistate Phase Change Materials”)) factoring algorithms (U.S. Pat. No. 6,714,954 (“Methods of Factoring and Modular Arithmetic”), modular arithmetic (U.S. Pat. No. 6,963,893 (“Methods of Factoring and Modular Arithmetic”)), and neural network processing (U.S. Pat. No. 6,999,953 (“Analog Neurons and Neurosynaptic Networks”). Applications of chalcogenide switching materials to data processing include U.S. Pat. No. 5,543,737 (“Logical Operation Circuit Employing Two-Terminal Chalcogenide Switches”).\nRecent work in the area of chalcogenide switching devices has demonstrated the operability of a three-terminal chalcogenide switching device. In these devices, a third terminal is added to the standard two-terminal chalcogenide switching device to enable control over the operating characteristics of the device. Application of a voltage signal or electric field to the third terminal, for example, provides a mechanism for controlling the magnitude of the threshold voltage needed to effect the switching transition between the other two terminals of the device. (U.S. Pat. Nos. 6,967,344 (“Multi-Terminal Chalcogenide Switching Devices”) and 6,969,867 (“Field Effect Chalcogenide Devices”, the disclosures of which are incorporated by reference herein).\nWith the advent of new chalcogenide devices having increased functionality, it is desirable to consider their potential to further expand the capabilities of chalcogenide-materials in the realm of computation. In particular, it is desirable to consider the suitability of three-terminal chalcogenide devices for applications in data storage or data processing and to devise device structures and circuits that exploit the capabilities of three-terminal devices. In U.S. Pat. Nos. 6,967,344 ('344 patent) and 6,969,867 ('867 patent), the disclosures of which are hereby incorporated by reference herein, Ovshinsky et al. further develop the notion of phase change computing by presenting additional computing and storage devices. The '344 patent discusses a multi-terminal phase change device where a control signal provided at one electrical terminal modulates the current, threshold voltage or signal transmitted between other electrical terminals through the injection of charge carriers. The '867 patent describes a related multi-terminal device that utilizes a field effect terminal to modulate the current, threshold voltage or signal transmitted between other terminals. The devices described in the '344 and '867 patents may be configured to provide a functionality related to that of a transistor.\nIn addition to new storage and processing devices, progress in the field of chalcogenide electronics would further benefit from the introduction of logic circuits based on chalcogenide materials that are capable of performing one or more logic functions. In particular, it is desirable to develop logic circuits based on chalcogenide memory and/or switching devices. The utilization of two-terminal chalcogenide switching devices in logic circuits has been discussed in U.S. Pat. Nos. 5,543,737 ('737 patent) and 5,694,054 ('054 patent); the disclosures of which are incorporated by reference herein. The potential for chalcogenide electronics further expand through the development of logic circuits that utilize the beneficial properties of the three-terminal family of chalcogenide devices."}
{"text": "Workstations with one or more robots are known from practice in various forms. As a rule, they are designed as fully automatic stations and they are equipped with protective partitions, which mechanically prevent the access of workers or switch off the industrial robot, for accident prevention for the workers.\nFurthermore, there are efforts to allow humans to cooperate or collaborate with industrial robots, especially tactile robots. This is called human-robot cooperation (HRC). Tactile articulated arm industrial robots suitable for this are known from, e.g., DE 10 2007 063 099 A1, DE 10 2007 014 023 A1 and DE 10 2007 028 758 B4.\nOn the other hand, a safety monitoring at an industrial robot by means of desired values is known from DE 10 2007 059 481 A1. DE 10 2005 027 522 B4 discloses an industrial robot with a moving monitored area around the tool or the dangerous area by means of a capacitive sensor system. A risk assessment according to DIN EN ISO 12100 is to be performed for designing workstations."}
{"text": "1. Field of the Invention\nThe present invention relates to a film combination or multilayer film with at least two layers, for particular use as a release film for membranes that contain oil.\n2. The Prior Art\nA large number of different film combinations have been disclosed. The problem with them, however, is either that the combinations do not have any barrier properties against oils, mineral oils in particular, or that they are not very stable, i.e. they are susceptible to tearing, have poor tear propagation properties or are not very thermally and/or mechanically stable."}
{"text": "1. Field of the Invention\nThe present invention relates to a drill bit sharpener, and more particularly to a drill bit sharpener having a rotating grinding wheel with an adjustable position to sharpen drill bits, which can be assembled quickly.\n2. Description of Related Art\nDrilling machines and drills make holes in objects. A drill has a drill bit that is turned to make the hole in the object. A cutting edge with a relief angle is formed on one end of the drill bit. After an extensive use, the cutting edge and the relief angle of the drill bit will become dull or break. A user must replace the dull or broken drill bit with a new drill bit to keep making holes efficiently, but replacement of the drill bits is not economical.\nTo sharpen a dull or broken drill bit is more economical than throwing the drill bit away. A conventional bench grinder with a rotating grinding wheel is often used to sharpen drill bits. The user holds the drill bit in his or her hands and uses the rotating grinding wheel to grind a cutting edge with the appropriate relief angle on the drill bit. Using a conventional bench grinder to sharpen drill bits has some shortcomings in both operating and using. The shortcomings include the following.\n1. Inconvenience operation:\nBench grinders are big and occupy a big space. The bench grinder must be mounted in a location before it can be used. The bench grinder is not easy to carry and not portable. Furthermore, holding the drill bit at the proper angle against the grinding wheel to accurately the tip is not easy. A user often much time to grind proper shape and angle on the drill bit.\n2. Hazardous operation:\nThe grinding wheel rotates at a very high speed, and at least part of the grinding wheel is exposed to surrounding environment. The user holds the drill bit in his or her hands and may inadvertently touch the grinding wheel if he.or she is the least bit careless. The high-speed, grinding wheel will injure the user if the user touches the wheel.\nTo overcome the shortcomings, the present invention provides a drill bit sharpener to mitigate or obviate the aforementioned problems.\nThe main objective of the invention is to provide a drill bit sharpener that will sharpen drill bits easily, quickly and efficiently.\nAnother objective of the invention is to provide a drill bit sharpener that is portable so a user can carry the sharpener and use it on the work site.\nOther objectives, advantages and novel features of the invention will become more apparent from the following detailed description when taken in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the variety designated CV297857, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include greater yield, better stalks, better roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen comes to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants and have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
{"text": "In recent years, attention has focused on methods for manufacturing electronic devices by ink-jet technology.\nAccording to a piezoelectric element in Patent Literature 1, electrodes provided on two ends along a short side of the piezoelectric element can be removed from one side of the piezoelectric element. The structure will be discussed below.\nFIG. 7(a) is a plan view of a first ceramic green sheet 30. FIG. 7(b) is a plan view of a second ceramic green sheet 40. Hereinafter, the first ceramic green sheet 30 will be called a first sheet and the second ceramic green sheet 40 will be called a second sheet.\nIn the conventional piezoelectric element, as shown in FIG. 8, the first sheets 30 and the second sheets 40 are stacked on a lower part 90, and then are pressed and burned thereon.\nA first electrode 50 of conductive paste is formed on the first sheet 30. A second electrode 60 is formed on the second sheet 40. The first electrode 50 has a first non-electrode region 51 and the second electrode 60 has a second non-electrode region 61.\nThe first and second sheets 30 and 40 are stacked, pressed, and burned so as to obtain a block 20 as an intermediate product shown in FIG. 8(a). FIG. 8(a) is a perspective view of the block 20. The front surface will be called a front surface 20a, the rear surface will be called a rear surface 20c, the right lateral surface will be called a right lateral surface 20b, and the left lateral surface will be called a left lateral surface 20d. \nThe first electrodes 50 are exposed only on the front surface 20a of the piezoelectric element 20 but are not exposed on the right lateral surface 20b, the left lateral surface 20d, and the rear surface 20c of the piezoelectric element 20.\nThe second electrodes 60 are exposed on the rear surface 20c but are not exposed on the front surface 20a, the right lateral surface 20b, and the left lateral surface 20d. Connection electrodes 70 are exposed on the front surface 20a, the right lateral surface 20b, and the rear surface 20c of the piezoelectric element. In other words, the connection electrodes 70 are formed up to the front surface 20a and the rear surface 20c of the piezoelectric element 20 that are opposed to each other.\nFIG. 8(b) shows the rear surface 20c of the piezoelectric element 20. The connection electrodes 70, the second electrodes 60, and the second non-electrode region 61 are shown on the front surface.\nAs shown in FIG. 8(c), a conductive layer 100 is formed over the rear surface 20c and requires more piezoelectric elements. As shown in FIG. 9, a plurality of slits 91 are formed so as to constitute the piezoelectric element 20.\nThe piezoelectric element 20 includes the connection electrodes 70 and the second electrodes 60 that are electrically connected to each other on the rear surface 20c via the conductive layer 100. Thus, the first electrodes 50 and the connection electrodes 70 are connected to external electrodes on the front surface 20a, activating a piezoelectric substance. Hence, external connection can be made only on one of the surfaces of the piezoelectric element 20. The surface is connected to a substrate having a number of external connection electrodes, thereby driving the piezoelectric elements. The piezoelectric element 20 used for an inkjet requires quite a number of piezoelectric substances, and a plurality of external electrodes need to be removed from one of the surfaces (front surface 20a) of the piezoelectric element 20.\nFIGS. 10(a) to 10(c) show displacements when a voltage is applied to the piezoelectric element 20. Specifically, an upper view shows a cross section of the piezoelectric element 20 while a lower view shows a displacement in this state. In this case, the cross section is perpendicular to the front surface 20a and the rear surface 20c of the piezoelectric element 20 in FIG. 8.\nIn FIG. 10(a), symmetric displacements appear when d1=d2 is satisfied where d2 and d1 are the widths of the second non-electrode region 61 and the first non-electrode region 51 on the two ends of the piezoelectric element 20.\nAs shown in FIG. 10(b), when “d1>d2” is satisfied, the displacement of the smaller d2 is larger than that of the larger d1. As shown in FIG. 10(c), when “d1<d2” is satisfied, the displacement of the smaller d1 is larger than that of the larger d2.\nAs described earlier, the locations of the maximum displacement of the piezoelectric element 20 vary depending upon the widths d1 and d2 of the non-electrode regions 51 and 61. In other words, the maximum displacement position and the maximum displacement of the piezoelectric element 20 vary depending upon the values of the widths of the non-electrode regions.\nThe values of the widths of the non-electrode regions vary because of the forming accuracy of the formed electrodes, a displacement when the ceramic green sheet is stacked, the degree of shrinkage during burning, and so on."}
{"text": "The CVD (chemical vapor deposition) method is one of methods of forming an insulating film in the process of LSI manufacture. According to the CVD method, various insulating films can be formed on a semiconductor substrate by changing the kind of material gas.\nTable 1 shows the relationship between material gases and insulating films (hereinafter referred to as CVD films) formed by means of the CVD method using the material gases.\nTABLE 1 ______________________________________ RELATION OF CVD FILM AND MATERIAL GASES CVD FILM MATERIAL GASES ______________________________________ SiO.sub.2 ##STR1## PSG MATERIAL OF SiO.sub.2 + PH.sub.3 BSG MATERIAL OF SiO.sub.2 + B.sub.2 H.sub.6 ASG MATERIAL OF SiO.sub.2 + AsH.sub.3 BPSG MATERIAL OF BSG AND PSG PO(OCH.sub.3).sub.3 GSG MATERIAL OF SiO.sub.2 + GeH.sub.4 Si.sub.3 N.sub.4 ##STR2## SiO.sub.x N.sub.y MATERIAL OF SiO.sub.2 + MATERIAL OF Si.sub.3 N.sub.4 Al.sub.2 O.sub.3 ##STR3## TiO.sub.2 ##STR4## Ta.sub.2 O.sub.5 ##STR5## ______________________________________\nFor LSI, SiO.sub.2, PSG, ASG, BPSG, GSG, Si.sub.3 N.sub.4 and SiOxNy are often used as interlayer insulating films and passivation films.\nWhen a gas having Si--H combination (for example, SiH.sub.4, SiH.sub.2 Cl.sub.2, or the like) is used as a material gas in forming one of those films, the Si--H combination is formed partly in the resulting CVD film (SiO.sub.2, SiOxNy, or the like).\nH (hydrogen) contained in the CVD film is trapped in the gate oxide of transistors, the interface between the gate oxide film and the semiconductor substrate, and so on, causing variations in the threshold of the transistors and reducing the life of hot carriers.\nIt is therefore required that the CVD film contain as little hydrogen as possible.\nAs the CVD methods, there are known atmospheric pressure CVD, low pressure CVD, plasma CVD, and so on.\nIn the atmospheric pressure CVD and the low pressure CVD, since thermal energy causes a chemical reaction to take place, CVD film deposition is generally performed at a high temperature of 700.degree. C. or more. In contrast, in the plasma CVD, CVD film deposition is performed by producing radicals of particles (atoms, molecules) in a plasma and causing a chemical reaction to take place between active particles, thus allowing the CVD deposition to take place at a low temperature of 250 to 400.degree. C.\nThat is to say, insulating films formed by the plasma CVD (hereinafter referred to as P-CVD films) are often used as insulating films (interlayer insulating films and passivation films) on a metal interconnection layer made of aluminum or the like.\nHowever, the CVD films (the P-CVD film in particular) have the following drawbacks:\na. Electron trap PA1 b. Moisture resistance\nH in CVD film is considered to reach the gate oxide film and tunnel oxide film (SiO.sub.2) of a transistor and cut the combination between Si and O in the gate oxide film and tunnel oxide film during the formation of that film or in the process subsequent to the film formation.\nAnd it is supposed that electrons are trapped at the combination cut sites and the trapped electrons have an adverse effect, such as variations in threshold, on transistors forming memory cells of a semiconductor memory (DRAM, SRAM, nonvolatile memory, or the like) and transistors making up logic.\nInsulating films (passivation films in particular) for use with LSI are required to have a property of blocking moisture. This is because, when the moisture reaches metal interconnection lines made of metal, such as aluminum, it will corrode them.\nWhen the plasma CVD is performed using, for example, SiH.sub.4, N.sub.2 O, N.sub.2 as material gases, an SiO.sub.2 film is formed. However, it is known that the SiO.sub.2 film is inferior in moisture resistance to an SiN film.\nOn the other hand, when the plasma CVD is performed using SiH.sub.4, NH.sub.3, N.sub.2 as material gases, an SiN film is formed. The SiN film is very good in moisture resistance. However, since the SiN film contains a large amount of H (Si--H combination), it is known that a large number of electron traps are produced in the gate insulating film and the like.\nThe present invention is made to solve the drawbacks and its object is, first, to find such an amount of Si--H combination in a CVD film fabricated by material gases containing a gas having Si--H combination as permits variations in the threshold of transistors due to electron traps in the gate oxide film (or the tunnel oxide film) to remain unchanged from those in the case where there is no CVD film fabricated by material gases including a gas having Si--H combination and apply the result to semiconductor devices, and, second, to eliminate threshold variations due to electron traps in a CVD film fabricated by material gases including a gas having Si--H combination and improve the moisture resistance."}
{"text": "This disclosure relates to the concentration of light, and, more particularly, the concentration of light using an optical element.\nTypically, light concentrators operate to receive light incident over a range of angles less than an acceptance angle at an aperture. The light is concentrated onto a region (e.g., on an absorber) with an area smaller than the area of the aperture. The ratio of the aperture area to the smaller area is known as the geometric concentration C. The laws of thermodynamics set a theoretical upper bound, known in the art as the “thermodynamic limit,” to the concentration for a given concentrator configuration.\nMany types of solar concentrators have been studied including reflective and refractive devices. Concentrators may be imaging or non imaging, and may be designed to correct for various types of optical aberration (spherical aberration, coma, astigmatism, chromatic aberration, etc.). For example, D. Lyndon-Bell, Monthly Notices of the Royal astronomical Society, vol. 334, pp. 787-796 (2002), describes an aplanatic concentrator featuring primary and secondary reflectors. However, the efficiency of such concentrators is limited by the obscuration of the primary reflector by the secondary reflector.\nOptical concentrators may be applied, for example, in the conversion of solar energy to electricity (or other form of energy). The power that a photovoltaic solar cell can produce is a function of the incident sunlight. A typical solar cell can utilize efficiently many times the un-concentrated incident sunlight in typical settings, provided that the temperature of the solar cell does not increase excessively. Therefore, an optical concentrator can be employed to concentrate sunlight onto a photovoltaic cell to improve the output of the photovoltaic cell. The output will increase with the concentration factor. At appreciable concentration factors, cooling may be required, since the efficiency of some photovoltaic cells may decrease rapidly with increasing temperatures.\nOptical concentrators may be applied in a variety of other applications including, for example, imaging, photography, concentration of light from sources such as lasers or light emitting diodes (LEDs), etc."}
{"text": "Businesses, governmental organizations and other entities are increasingly saving large volumes of data necessary for daily operations. This data represents a significant asset for these entities. Consequently, data loss, whether accidental or caused by malicious activity, can be costly in terms of wasted manpower, loss of goodwill from customers, loss of time and potential legal liability. To ensure proper protection of data for business and legal purposes (e.g., to ensure quick recovery of data in the event of a disaster, to comply with document retention requirements, etc.), these entities often back up data to a physical media, such as magnetic tapes or optical disks on a regular basis.\nTraditional backup systems placed an application server, backup server, source device, destination device and a local area network (“LAN”) in the data path of backup operations. Under these systems, the LANs were becoming overburdened by the amount of data being copied. Often, the backup window (the period in which data unavailable for normal operations in order to permit backup) was too short to achieve a complete backup of data. Accordingly, many entities implemented Storage Area Networks (“SAN”) to relieve the burden of mass data storage and backup from the LAN, freeing the LAN for more immediate data storage and manipulation operations. In SANs, data from multiple machines on a network may be backed up to a remote media library. Centralized data backup allows storage problems to be identified at one location and has the advantage of increased efficiency.\nOne example of a media library commonly used in enterprise backup systems is a magnetic tape library. In a typical magnetic tape library, tapes are contained in cartridges and the tape library contains multiple cartridge slots in which tape cartridges can be stored. The tape cartridges are physically moved between cartridge slots and tape drives by a robot. The robot is controlled by access commands received from the host devices on the network. When specific data is required, the host device determines which cartridge slot contains the tape cartridge that holds the desired data. The host device then transmits a move-element command (e.g., the SCSI “move medium” command) to the robot and the robot moves the tape cartridge, mounting the cartridge in the appropriate tape drive.\nThe Linear or Linear Tape File System (LTFS) Format Specification by IBM and the LTO Program (hereby fully incorporated by reference in its entirety for all purposes) defines a file system for LTO-5 tapes, LTO-6 tapes and may be extended to other tapes using an eXtensible Markup Language (XML) schema architecture. This file system support allows the use of an LTFS-formatted tape as if it were a file system. Files and directories may appear in a directory listing, files may be dragged and dropped from tape, data may be accessed at the file level, etc.\nConsequently, while it previously was necessary to make use of a backup application to write and read tapes, the introduction of LTFS has simplified the storing and retrieval of files on tape by reducing such operations to a copy. Furthermore, any operating system that includes LTFS support can mount LTFS for a properly formatted tape and read and write the files thereon.\nBefore an LTFS tape can be used, it must be mounted in an LTFS-compatible tape drive, and the LTFS filesystem must be mounted. In this disclosure, the terms “file system” and “filesystem” may be used interchangeably. Once the LTFS filesystem has been mounted, files on the tape can be accessed in the same manner as files in other filesystems (e.g., written, read, dragged and dropped, etc.). When the tape is no longer needed, the LTFS filesystem is unmounted, and then the tape can be removed from the tape drive.\nIn some systems, the LTFS-compatible tape drives are accessed using Unix/Linux commands. The mounting and unmounting of the LTFS filesystem in these systems are asynchronous operations. In other words, when a mount or unmount command is issued, the operation is initiated, but the user does know when the operation is complete—no notification of the completion is provided by the operating system, and the operation does not complete with known timing. If a command which is dependent upon the completion of the mount or unmount operation is issued before the mount or unmount operation is completed, the subsequently issued command will fail. This problem is normally resolved by simply retrying the subsequent command. After some number of retries, the mount or unmount operation will presumably have had sufficient time to complete, and the subsequent command will be successful.\nWhen the mount/unmount and subsequent commands are being manually issued by a user, the solution of retrying the subsequent commands may be sufficient. In the case of automated systems (e.g., media library interface/management systems) that issue these commands, however, there are some problems with this solution. For example, when an automated, synchronous system retries the subsequent command before the preceding mount/unmount command is complete, an error message may be generated—repeated retries may result in a large number of error messages. A user may interpret the large number of error messages as an indication of a malfunction in the system, rather than a normal and expected result of the retries, and may consequently feel that it is necessary to interrupt operation of the system to determine the source of the errors. Additionally, the retries increase the resource usage of the system, and the delays between completion of the asynchronous commands and the succeeding retries reduce the efficiency of the system.\nIt would therefore be desirable to provide a mechanism for enabling a synchronous media library interface/management system to perform or interface with asynchronous mount and unmount operations in a way that eliminates the need to continually retry commands that follow these operations and increases the efficiency of the system."}
{"text": "The present invention relates to assembling teams to solve a problem, and more particularly relates to a system and method of assembling a team for solving a problem within an immersive environment that accounts for individual and group biases.\nThe selection of a team to solve a problem is fraught with any number of possible problems. Individuals and groups are prone to all types of biases that can lead a team to an incorrect conclusion. Moreover, team based outcomes can depend on the types of environments the members are immersed within.\nImmersive environments, social decision making, and the interactions between humans and systems produce data for the development of predictive and prescriptive models. At intersection are analytics, social psychology, and immersive technologies. The use of immersive environments within the context of social interaction paradigms enables the wisdom of many minds, e.g., following the Condorcet jury theorem. Condorcet's jury theorem is a theorem about the relative probability of a given group of individuals arriving at a correct decision.\nThe prevention of knowledge cocoons, information cascades, decision polarization, and the attachment to premature hypothesis is a core component of data analysis. The accelerating data avalanche where data is producing even more data perpetuates the need of turning information into knowledge for the consumption of parallel analysis. Such group, crowd, or team work sourcing produces precursors for the aggregation of fuzzy social decisions. Social or group dynamics produce fuzzy decisions or effective errors that are magnified as higher ply groups aggregate lower ply groups' decision until a root analysis is completed. Predictive social decision analytics reduces team effective errors by recommending a team, social paradigm, and immersive environment to maximize individual analysis for the correct team decision."}
{"text": "1. Field of the Invention\nThe present invention relates to semiconductor devices and, in particular, relates to DRAM devices having improved charge storage capabilities.\n2. Description of the Related Art\nThe Dynamic Random Access Memory (DRAM) device is an important component of electronic systems requiring digital storage capabilities, such as digital computers, personal data assistants (PDA), and the like. The typical DRAM provides an array of memory cells disposed in a high density configuration. For example, a single DRAM device having a minimum feature dimension of 0.15 micron can include an array of millions of directly addressable memory cells that are each capable of storing a single bit of digital data.\nFIG. 1 schematically illustrates a typical memory cell 20 comprising a single capacitor 21 and a field effect transistor (FET) 22. The capacitor 21 stores a quantity of charge which is indicative of the state of the memory cell 20. The FET 22 acts as a switch that provides a conducting path to the capacitor 21 when the cell 20 is being addressed. The FET 22 also isolates the capacitor 21 when the memory cell 20 is in a quiescent state, i.e., not being addressed, so that the capacitor 20 can store charge for extended periods of time.\nAs shown in FIG. 1, the FET 22 is formed within a substrate 23 and the capacitor 21 formed adjacent the substrate 23 in an ILD layer 31. The substrate 23 includes first and second doped regions 24, 25 and a channel 26 extending therebetween. For example, the first doped 24 region can be configured to be a source input of the FET 22 and the second doped region 25 can be configured to be a drain input of the FET 22. The FET 22 further comprises a gate electrode 27 disposed over the channel 26 which extends between the source 24 and drain 25 of the FET 22 and is insulated from the substrate 23 by a thin insulating layer 28. Furthermore, the capacitor 21 comprises first and second plates 29, 30 separated by an insulating layer 31, wherein the first plate 29 is coupled to the source of the FET 22 and the second plate 30 is coupled to a fixed reference voltage.\nAs is well known in the art, when the voltage applied at the gate 27 is greater than a threshold value, the channel 26 becomes a conducting path that allows charge to readily flow between the source 24 and drain 25. Thus, in response to a write signal arriving at the drain via a digit line 32, charge is able to flow through the channel 26 to the first plate 29 of the capacitor 21, thereby providing the capacitor 21 with a quantity of charge which is indicative of the input write signal. Likewise, during a read cycle, an output read signal can be developed at the drain 25 of the FET 22 which is indicative of the charge stored in the capacitor 21. Furthermore, during quiescent periods when the memory cell 20 is not being addressed, the gate voltage is reduced below the threshold value so as to inhibit conduction across the channel 26 and, thereby, help preserve the charge of the capacitor 21.\nA common problem with memory cells of the prior art is that the capacitor is unable to store charge indefinitely. For example, even if the gate voltage is maintained below the threshold value, a subthreshold leakage current usually flows through the channel 26 which discharges the capacitor during the quiescent period. To accommodate such discharging, typical DRAM devices also include refresh circuitry which periodically monitors the state of each memory cell and repeatedly recharges individual capacitors having a detected residual charge. However, because the DRAM is inaccessible during the refresh cycle, such refreshing reduces the rate at which data can be read from or written to the device. Thus, there is a need to reduce the subthreshold leakage currents so as to extend the time between refresh cycles and, therefore increase the throughput of the device.\nOne method that can be used to reduce subthreshold leakage currents involves disposing complementary dopant atoms in the channel region of each FET of the memory array. For example, if the source and drain of the FET include pentavalent dopant atoms, such as Phosphorous, the complementary dopant atoms would comprise trivalent dopant atoms, such as Boron. The presence of the complementary dopant atoms within each channel region increases the corresponding threshold gate voltage and, as a result, also increases the resistance of the channel region during quiescent periods.\nTo ensure that each FET of the memory array is configured with substantially uniform operating characteristics, i.e. uniform gate threshold values, it is important for each channel region to be doped in a substantially identical manner. For example, a first channel region having complementary dopant atoms disposed throughout the first channel region will provide a substantially different gate threshold value than that of a second channel region having a reduced concentration of complementary dopant atoms near its edges. Thus, because of the difficulty of precisely embedding complementary dopant atoms only within the relatively narrow confines of the channel of the typical memory array using conventional masking techniques, the industry has adopted the practice of doping the complementary atoms using a blanket implant process. In other words, no attempt is made to prevent some of the complementary atoms from becoming implanted within regions of the substrate immediately adjacent the channel regions.\nThe blanket implant process of the prior art comprises exposing a relatively large portion of the substrate corresponding to the entire memory array to a source of energetic complimentary dopant atoms. As a result, complementary dopant atoms 33 are embedded not only within the channel 26 but also in the doped source and drain regions 24, 25 that surround the channel 26 as shown in FIG. 1 such that the concentration of complementary dopant atoms is substantially uniform across the source 24, channel 26 and drain 25. Unfortunately, the increased concentration of complementary dopant atoms 33 in the source region 24 effectively forms a diode junction with adjacent regions of the substrate 23. Consequently, the blanket implant process of the prior art produces a relatively large junction leakage current that flows from the source 24 directly into the surrounding substrate 23. Thus, the reduction in the subthreshold leakage current is gained at the expense of the increase in the junction leakage current which is undesirable since it contributes to the discharging of the capacitor 21.\nFrom the foregoing, therefore, it will be appreciated that there is a need for a memory array having reduced discharge rates and methods for providing the same. In particular, to provide reduced subthreshold leakage currents, there is a need for the channel regions of the substrate corresponding to the array of memory cells to include complementary dopant atoms disposed so as to increase the threshold gate voltage. Furthermore, to reduce junction leakage currents, there is a need for the complementary dopant atoms to be substantially absent from actively doped regions of the substrate that are coupled to the charge storing capacitors. Finally, there is a need for the complementary dopant atoms to be disposed in such a way that the channel regions of the memory arrays becomes conductive at substantially uniform threshold voltages.\nThe aforementioned needs are satisfied by the memory array and method of manufacturing the same of the present invention. In one aspect, the present invention comprises a memory array of a DRAM device that comprises a plurality of memory cells for storing information, wherein the memory cells are arranged into a plurality of cell pairs such that each pair includes first and second cells each comprising a storage element for storing charge and a valve element for controlling the flow of charge to and from the corresponding storage element. The valve elements of each cell pair are disposed adjacent each other and the charge storage element of the cell pair are outwardly disposed form the valve element of the cell pair. The array further comprises a substrate comprising a plurality of valve regions and a plurality of storage regions wherein the first and second valve elements of each cell pair overlap a corresponding valve region of the plurality of valve regions. The first and second storage elements of each cell pair are electrically coupled to respective first and second storage regions of the plurality of storage regions and the valve regions of the substrate include a first concentration of complementary dopant atoms. The storage regions of the substrate include a reduced second concentration of complementary dopant atoms, thereby reducing the rate at which charge escapes from the charge storage element during quiescent periods.\nIn another aspect, the present invention comprises a method of reducing leakage currents in a DRAM device. The method comprises selectively implanting complementary dopant atoms within a first region of the substrate and then forming a plurality of field effect transistors each comprising first and second active nodes and a channel extending therebetween. The first active node is disposed in the substrate so that a substantial portion of the first active node is disposed outside of the first region so as to reduce a junction leakage current and the channel and second active node are disposed within the first region so as to reduce subthreshold leakage current."}
{"text": "1. Field of the Invention\nThe present invention relates to a call device, a call control system, a call management system, and a call control method.\n2. Description of the Related Art\nThere have heretofore been available in the art cellular phone terminals which output a melody (hereinafter referred to as “ringer melody”) that is preset by the user of a cellular phone terminal as a recipient's terminal during a period of time after the cellular phone terminal has received an incoming call until the user starts to talk to the calling party. There have also been available in the art cellular phone terminals which output a melody (hereinafter referred to as “ring back melody”) that is preset by the user of a cellular phone terminal for a recipient's terminal during a period of time after a cellular phone terminal has sent an outgoing call until the parties start talking to each other.\nThe user of a cellular phone terminal presets a preferred melody as a ringer melody or a ring back melody. Therefore, the user, who owns a cellular phone terminal which outputs a ringer melody or a ring back melody, can enjoy the music until the parties start talking to each other through the cellular phone terminal.\nThere has been proposed a call device, which is capable of superimposing a music signal on a speech signal representative of speech for the user to enjoy the music while talking to the other party.\nFor example, Document 1 (Japanese laid-open patent publication No. 2002-101206) discloses a telephone communication service management device for superimposing speech advertisement information such as music information on speech information and outputting the combined information. The telephone communication service management device presets the speech advertisement information.\nDocument 2 (Japanese laid-open patent publication No. 2002-344571) reveals a call device for superimposing a music signal on a speech signal and outputting the superimposed signal. The user of a cellular phone terminal presets the music signal.\nRinger melodies or ring back melodies give a certain level of enjoyment to the users of cellular phone terminals.\nHowever, it is desirable for ringer melodies or ring back melodies to give further enjoyment to the users of cellular phone terminals.\nAt present, a cellular phone terminal outputs a ringer melody or a ring back melody for a short period of time from the arrival of an incoming call to the start of the connection. Therefore, the user of a cellular phone terminal can enjoy the music that is preset by the user of a cellular phone terminal as a recipient's terminal, for a short period of time from the arrival of an incoming call to the start of the connection.\nAccording to Document 1, the music and the speech which are outputted together, are provided by the telephone communication service management device, without any reference to the called party. Therefore, the telephone communication service management device may possibly output music together with speech, which the called party is not fond of.\nThe telephone communication service management device disclosed in Document 1 is also disadvantageous in that if a ringer melody or a ring back melody is preset, then the music, which is output after the arrival of an incoming call until the parties start to talk to each other, is different from the music, which is output while the parties are talking to each other. Therefore, the user of a cellular phone terminal can enjoy the music (a ringer melody or a ring back melody) that is preset by the user of a cellular phone terminal for a recipient's terminal, only during a short period of time after the arrival of an incoming call until the connection is established.\nWith the call device disclosed in Document 2, if a ringer melody or a ring back melody is preset, then the music, which is output from the arrival of an incoming call until the parties start to talk to each other, may possibly be different from the music, which is output while the parties are talking to each other. Consequently, the user of a cellular phone terminal can enjoy a ringer melody or a ring back melody, only during a short period of time after the arrival of an incoming call until the connection is established.\nAnother problem is that the called party who has preset a ring back melody is unable to confirm the amount of the ring back melody to which the calling party has already listened. Therefore, the called party cannot judge whether the calling party has fully enjoyed the ring back melody or not."}
{"text": "The present invention relates to a control circuit for controlling at least one solenoid valve for fuel metering in an internal combustion engine.\nSuch a control circuit is for example discussed in German Published Patent Application No. 195 39 071. Using this control circuit, rapid solenoid valves are controlled, e.g., for common rail fuel injection or gasoline direct injection of internal combustion engines, via booster and battery supply voltage FET switches, and the energy is released during the transition from the pickup current to the retaining current being stored in a capacitor.\nThis circuit system may require a multiplicity of components and complex driver circuits that make it possible to use large currents to drive the valves to be controlled. Here, the resulting dependencies of the cutoff edges during the rapid disconnection of the solenoid valves from the battery voltage may be disadvantageous.\nAn object of an exemplary embodiment of the invention is to provide a control circuit for controlling at least one solenoid valve for fuel metering in an internal combustion engine, so that the cutoff edges during the rapid disconnection are largely independent of modifications of the battery voltage.\nThe control circuit according to an exemplary embodiment of the present invention includes a recharge circuit that is connected with the first and second terminal of the supply voltage to produce a pre-stabilized recharge voltage for the storage capacitor from the battery supply voltage, and is connected with the first switching apparatus, arrangement or structure, the recharge circuit supplying, via third switching apparatus, arrangement or structure, situated between the first terminal of the at least one solenoid valve and the recharge circuit, the solenoid valves with current from the energy stored in the storage capacitor in the booster phase, the recharge circuit also containing the storage capacitor as well as a fourth switching apparatus, arrangement or structure, that, controlled by the control apparatus, arrangement or structure, activates the recharge circuit in order to recharge the storage capacitor.\nThe dependence of the solenoid valve shutoff on modifications in the battery voltage is omitted to the greatest possible extent. The usable range of operation is thus expanded.\nIt is believed that the battery voltage can include a large range without influencing the shutoff (or operating, clearing) time of the solenoid valves. Here, the stabilized voltage can lie above or below the battery voltage."}
{"text": "The present disclosure relates to a continuous-time system delta-sigma (ΔΣ) modulator and a signal processing system which are applied to audio equipment and medical instruments in addition to receivers in wireless communication.\nReferring to FIG. 1A and FIG. 1B, there are shown conceptual diagrams illustrating general continuous-time system quadrature lowpass type ΔΣ modulators. FIG. 1A shows a feedback-type ΔΣ modulator and FIG. 1B shows a feedforward-type ΔΣ modulator.\nA ΔΣ modulator 1A shown in FIG. 1A is made up of integrators INT1, INT2, INT3, and INT4, a quantizer Quan, adders ADD1, ADD2, ADD3, and ADD4, and digital-to-analog (DA) converts DAC1, DAC2, DAC3, and DAC4.\nA ΔΣ modulator 1B shown in FIG. 1B is made up of integrators INT1, INT2, INT3, and INT4, a quantizer Quan, adders ADD1 and ADD2, and a DA converter DAC1.\nIn FIGS. 1A and 1B, u is representative of an analog input signal and v is representative of a digital output signal. a1, a2 and a3 are each representative of a feedback gain or a gain in feedforward path. Q is representative of an effective gain of the quantizer. The gain of the DA converts DAC4 and DAC1 that give feedback to the integrator INT4 is 1.\nA noise transmission function (NTF) to digital output signal v of quantization noise generated in the quantizer Quan is expressed by equation (1) below in both FIG. 1A and FIG. 1B, indicating a frequency characteristic of highpass type.\n                              N          ⁢                                          ⁢          T          ⁢                                          ⁢          F                =                              s            4                                              s              4                        +                                          a                1                            ⁢                              Qs                3                                      +                                          a                2                            ⁢                              Qs                2                                      +                                          a                3                            ⁢              Qs                        +            Q                                              (        1        )            \nTo be more specific, a quantization noise generated in the quantizer by a feedback effect in the ΔΣ modulator is noise-shaped to be shifted to high-frequency area, thereby providing a high S/N ratio within a signal band.\nReferring to FIG. 2, there is shown a circuit diagram illustrating an exemplary configuration of an integrator providing an element block of the ΔΣ modulator.\nReferring to FIG. 3, there is shown a circuit diagram illustrating an exemplary configuration of an adder providing an element block of the ΔΣ modulator.\nThe integrator INT shown in FIG. 2 has an OTA (Operational Transconductance Amplifier) 1, a resistor element R, and a capacitor C.\nThe adder ADD shown in FIG. 3 has an OTA2, and a resistor elements R1 and R2.\nA transfer characteristic between input and output of the integrator INT is expressed as equation (2) below by use of a resistor (R) and a capacitor (C).\n                              T          ⁡                      (            s            )                          =                                            v              out                                      v              in                                =                      1            sCR                                              (        2        )            \nLikewise, a transfer characteristic between input and output of the adder ADD is expressed as equation (3) below by use of two resistors R1 and R2.\n                              T          ⁡                      (            s            )                          =                                            v              out                                      v              in                                =                                    R              2                                      R              1                                                          (        3        )            \nAs shown in FIG. 2 and FIG. 3, the integrator INT and the adder ADD can each be configured by use of one amplifier OTA.\nOnly if the adder ADD is arranged immediately before the integrator INT, one amplifier can realize both an integrator and an adder.\nReferring to FIG. 4, there is shown a circuit diagram illustrating a circuit realized by an adder and an integrator having two inputs and one output.\nOutput voltage vout of this circuit can be expressed as equation (4) below, providing an operation realized by an adder and an integrator combined.\n                                              ⁢                              V            out                    =                                                    1                                  sCR                  1                                            ⁢                              V                                  in                  ⁢                                                                          ⁢                  1                                                      +                                          1                                  sCR                  2                                            ⁢                              V                                  in                  ⁢                                                                          ⁢                  2                                                                                        (        4        )            \nNow, paying attention to the ΔΣ modulator shown in FIG. 1, it is seen that the modulator can be configured by four amplifiers because all adders are arranged immediately before the integrator in the case of feedback type shown in FIG. 1A.\nOn the other hand, in the case of the feedforward type shown in FIG. 1B, the adder arranged immediately before the quantizer Quan cannot form a circuit with the integrator, so that five amplifiers are required, thereby leading to increased power and increased circuit scale.\nIn order to avoid a situation where one more amplified is required, a ΔΣ modulator of feedforward type that does not use the adder arranged immediately before the quantizer may be considered.\nReferring to FIG. 5, there is shown a configuration of a ΔΣ modulator that does not use any adder immediately before a quantizer as a variation to the configuration shown in FIG. 1B.\nIn a ΔΣ modulator 1C shown in FIG. 5, every add operation is executed in front of an integrator INT1 of the last stage, so that every feedback signal passes two integrators.\nPassing two integrators raises the degree by two degrees, so that the feedforward path including a1 is set to sa1 rather than a constant, thereby realizing first-degree coefficient by providing a differential characteristic.\nReferring to FIG. 6, there is shown a circuit for realizing a portion enclosed by dotted lines shown in FIG. 5.\nA feedforward path from the output of the integrator INT4 to the input of the integrator INT1 can be realized by coupling, by use of a capacitor, the input and the output of an amplifier OTA configuring a corresponding integrator."}
{"text": "The use of conventional wind chimes to produce pleasing sounds in response to air movement is very well known. Although such wind chimes exist in many different structural configurations, they typically include one or more suspended vanes. In response to sufficient air motion, the vanes collide against one another and/or against a proximate housing, and produce sound determined primarily by the natural frequency of the colliding members."}
{"text": "This invention relates to an electrical contact assembly used for a circuit breaker of a control center.\nIn a known control center including circuit breakers, control units constructed by control means, such as a circuit breaker and an electromagnetic contactor, are accommodated in multistage in a frame board of the control center so that the control units can be drawn out therefrom. In the circuit breaker, a movable contact is arranged to contact with or separate from a stationary contact disposed at the rear portion of the frame board in accordance with the movement of the control unit in a direction to draw out the movable contact from the frame board.\nThe circuit breaker unit includes interrupting units located on a power source side at which a bus bar and a control unit are interconnected and on a load side connected to an external cable of the control unit.\nFIG. 1 is a general perspective view showing a control unit located on the load side, in which reference numeral 3 designates a control unit including a distribution circuit breaker 1 and an electromagnetic contactor 2, and bar conductors 5 are used as a movable contact unit 4a on the movable side. FIG. 2 shows a frame board 6 for accommodating the control units 3 in multistage and a plurality of compartments for accommodating control units are formed by inserting plates 7 into the frame board 6 at different levels. On the side of each compartment 8 is attached a stationay contact unit 4b which is connected to or separated from the movable contact unit 4a shown in FIG. 1.\nFIGS. 3a and 3b show a general structure of a stationary main circuit breaker unit, in which U-shaped contacts 11 of the stationary contact unit 4b are provided to contact with or separate from the bar conductors 5 on the movable side 4a. Each of the U-shaped contacts 11 is prepared by forming a plate conductor into a U-shape having contacting portions 11a which are formed by slightly bending the sides thereof inwardly to form a narrow opening for clamping the movable contact therebetween. Although the bar conductor 5 is caused to engage and clamped by the contacting portions 11a, the U-shaped contact 11 can be made of a resilient electrically conductive material or the bar conductor may be sandwiched by using another suitable resilient members, not shown. The U-shaped contacts 11 are supported by a holder 12 provided with recesses 13 for receiving the contacts 11.\nFIG. 4 shows a cross-sectional view taken along the line IV--IV in FIG. 3b, in which the bottom of the U-shaped contact 11 is connected to one end of a connecting conductor 14 and to the other end thereof is secured a screw 15 which connects the conductor 14 to an external cable, not shown.\nWhen the control unit 3 is inserted into the compartment 8, the bar conductor 5 is also inserted between the narrowed contacting portions 11a of the U-shaped contact. In this inserted condition, current flows through the contacts and if a large current accidentally flows in such a case as short circuit, a large electromagnetic repulsive force is generated in proportions to the current flown, thereby widely separating the opening between the contacting portions 11a thus generating electric arc which damages the U-shaped contact 11. Therefore, in order to avoid such separation phenomenon or damage described above, it is desired to preliminarily impart contact force to the contacting portions 11a to overcome the repulsive force generated at a time when large current flows. However, in a case where a U-shaped contact with relatively large contact force is used, a large operating power is required to insert the control unit 3 into and draw out it from the frame board 6. In addition, since the connecting conductor 14 is sometimes bent by strong vibration of the cable caused by the electromagnetic force due to short circuit current, it is necessary to detachably connect the conductor 14.\nFIG. 5 shows another example of a prior art circuit breaker unit on the movable side, in which a U-shaped contact 32 is arranged to contact with a stationary bar conductor 31, for example, a bus bar. In this example, the U-shaped contact 32 is formed by using two resilient conductors and the opening thereof is narrowed so as to form contacting portions 32a between which the bar conductor 31 is to be inserted and clamped.\nThe U-shaped contacts 32 are supported by a holder 33 covered by a cover 34 and provided with recesses 35 for receiving the U-shaped contacts 32. A projection 36 is formed in the recess 35 for guiding the contact 32. FIG. 6 shows a cross-sectional view taken along the line VI--VI in FIG. 5, in which the U-shaped contact 32 is connected to an electric cable 41 by, for example, welding. When the movable contact unit is moved to engage the U-shaped contact 32 with the bar conductor 31, the projection 36 serves to guide the contact 32 so as to cause the center line thereof to coincide with that of the bar conductor 31, and the U-shaped contact 32 is separated from the bar conductor 31 by moving the contact unit leftwardly in FIG. 6.\nWith the example shown in FIGS. 5 and 6, the same adverse result as described with reference to the former example shown in FIGS. 1 through 4 is caused when a large current accidentally flows. In addition, regarding the latter example, it is difficult to sufficiently utilize the space in the control unit for the reason that the cable 41 connected to the contact 32 has a vertical portion 40 as shown in FIG. 6 and another member cannot be disposed thereon."}
{"text": "The invention relates to a lost foam casting system, and more particularly to an improved method and apparatus for introducing molten metal to the pattern assembly.\nLost foam casting is a known technique. A pattern assembly formed of evaporative foam material, for example gasifiable or liquifiable expanded polystyrene, or expanded polymethylmethacrylate (PMMA), either of which which may be coated with a thin ceramic coating, is placed in a flask and surrounded by unbonded particulate media, most commonly sand, though other particles can be used such as zirconia, metal spheres, etc. Upon application of molten metal to the pattern assembly, the foam material vaporizes and is replaced by the metal in the shape of the pattern assembly The vaporized foam material escapes into the interstices in the sand.\nThe pattern assembly typically includes one or more workpieces of given respective patterns for yielding cast metal parts, and a sprue connected by one or more in-gates to the workpieces for communicating the molten metal through the sprue to the in-gates and then to the workpieces. After cooling, the cast metal sprue and the in-gates are broken away from the workpieces and discarded or recycled. The cast metal workpieces are then trimmed to yield the final product.\nThe present invention arose out of efforts to provided improved cast aluminum alloy cylinder heads and engine blocks for internal combustion engines in marine applications, and cast steel two cyle outboard crankshafts, though the invention is not limited thereto.\nIn the present invention, molten metal is applied to the foam pattern assembly by vacuum lift. Vacuum is applied through the sand, and a gas-permeable member or path is provided communicating between the sand and a vertical fill passage such that molten metal is vacuum lifted to the foam pattern assembly."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of connectors.\n2. Prior Art\nU.S. Pat. No. 5,533,915 discloses a connector suitable for use in relatively high current applications such a encountered in certain battery operated devices. The connector is a two conductor connector configured to avoid the possibility of being plugged in backward, and having beryllium copper spring members on each male conductor to force the male conductor against the female conductor on insertion of the male connector part into the female connector part."}
{"text": "1. Field of the Invention\nThe present invention relates to novel azo compounds and a process for preparing the same and, more particularly, to the azo compounds and a process for preparing the azo compounds represented by the formula: ##STR2## wherein X is hydrogen or halogen. The azo compounds are useful as precursors of disperse dyes and basic dyes and are effective for therapeutic treatment.\n2. Brief Description of Prior Art\nGenerally, azo compounds are known to be useful as dyes including disperse and basic dyes. It is also known that imidazole can be a component to form various azo compounds. It is not known, however, that the azo compounds of the above formula can be basic and disperse dyes and that in particular they are effective for therapeutic treatment.\nIt is generally known that, when imidazole is used as a coupling component of azo compounds, it is usually employed in an amount of from about 1 to 1.5 moles based on a diazo component to which imidazole is coupled. In this case, however, the reaction permits the production of by-products in which the diazo component may be coupled to the 4-imidazolyl group or one more diazo components may be coupled to the 4-imidazolyl group. Accordingly, this reaction has a drawback in not providing the azo compounds in a pure form."}
{"text": "The present disclosure relates generally to the field of amusement parks. More specifically, embodiments of the present disclosure relate to an interactive component for an amusement park.\nThis section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present disclosure, which are described below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present disclosure. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nAmusement parks contain a variety of rides providing unique experiences to each park guest. The addition of large attractions, such as rides and shows, generally provides an amusement park with additional capacity to handle a larger number of guests. However, such attractions tend to draw more visitors and become hubs for guest traffic. Further, the addition of traditional rides without an added layer of intrigue may be insufficient to garner sufficient guest interest to address either guest traffic issues or provide an advantage over competitors. In some cases, amusement parks may include a variety of interactive features that are utilized to entertain park guests when waiting for and/or travelling between the large attractions. With the increasing sophistication and complexity of modern attractions, and the corresponding increase in expectations among amusement park and/or theme park guests, improved and more creative attractions are needed, including attractions that provide a more interactive and personalized experience."}
{"text": "It is believed that a signal receiver and a signal transmitter for wirelessly receiving and transmitting signals, respectively, will eventually become standard equipments of a computer in view of the current trend that the data communications between two computers are brought about by remote data signals, and that more and more computer users are on the buying spree of such a signal receiver and such a signal transmitter, which are hooked up with the computer and its peripheral equipments.\nIt is rather conceivable that such a practice, as described above, of adding the signal receiver and the signal transmitter for wirelessly receiving and transmitting signals, respectively, to a computer system is not only expensive but also wasteful."}
{"text": "In a paper conveying apparatus of an image reading apparatus, image copying apparatus, etc., sometimes a jam occurs when the paper moves along the conveyance path. In general, a paper conveying apparatus is provided with the function of determining whether a jam has occurred by a paper being conveyed to a predetermined position inside the conveyance path within a predetermined time from the start of conveyance of the paper and of stopping the operation of the apparatus when a jam has occurred.\nOn the other hand, if a jam occurs, a large sound is generated in the conveyance path, so the paper conveying apparatus can determine whether a jam has occurred based on the sound which is generated on the conveyance path and thereby detect the occurrence of a jam without waiting for the elapse of the predetermined time.\nA jam detection apparatus of a copier which converts the sound which is generated on the conveyance path to an electrical signal and determines that a jam has occurred when the time when the signal is over a reference level exceeds a reference value has been disclosed (see Japanese Laid-open Patent Publication No. 57-169767)."}
{"text": "The proliferation of mobile communication devices including smart phones, notepads and laptops has resulted in significant efforts for efficient wireless transmission through adoption of new standards, hardware improvements and more sophisticated network management. As a result, there has been an ever increasing pressure to reduce cost, size and power consumption of such systems. This has driven designers to develop transceivers with higher levels of integration with much of the work focused on replacing the external components with integrated counterparts. This is not a trivial matter and often necessitates design of completely new transceiver architectures to achieve fewer off chip components.\nOn the other hand, there is a very high demand for performance from the modern integrated transceivers in order to achieve high modulation orders and hence to fulfill the high data rates driven by today's wireless and mobile applications. In this respect, the transceivers are required to provide an exceptionally good noise figure performance, linearity, matching, IQ (in-phase/quadrature) phase/gain balance and radio frequency (RF) local oscillator (LO) phase noise performance. While the latter is determined by the performance of the synthesizer circuit, the rest of the parameters are dependent on the components used within the transmit and receive chains.\nGiven the emphasis on integration and low cost, there is a major push for the adoption of, for example, sub-micron CMOS technology for the implementation of transceiver technology. The key characteristic of the technology is the low cost for implementing the digital functions but poor analogue and RF performance compared with other specialized technologies. Hence any use of digital functions in CMOS for improving the transceiver performance is very desirable and enables the use of the technology in high data rate wireless communication systems."}
{"text": "1. Field of the Invention\nThis invention relates to a receptacle holder for mounting along a horizontal support. More particularly, the invention relates to a device enabling to mount any kind of receptacles, such as garbage pails, and supports of all kinds on the outside of the railing of a veranda.\n2. Description of the Prior Art\nIt is well known that small residential verandas very rapidly become clogged with a variety of articles, such as tables, chairs, various tools and garbage pails or bags, with the result that there is a very limited area for enjoying life outside. On the other hand, it is traditional especially in multi-dwelling houses, to place garbage pails or bags on the veranda, since this is the most appropriate place.\nThere is therefore a need to dispose of some of the material which is normally found on the veranda while making the additional space readily available to the user.\nSome receptacle holders are known in the prior art, such as those found in the following United States Patents:\n1,072,022 PA1 2,550,019 PA1 2,908,386 PA1 3,007,664 PA1 3,016,149 PA1 3,892,315 PA1 3,980,202 PA1 4,643,380\nHowever, either in view of their cost or their construction, these holders could not be easily adapted for mounting on the outside of the railing of a veranda."}
{"text": "1. Field of the Invention\nThis invention relates to a composition suitable for use in mechanical power transmission units. More particularly, the invention relates to compositions specifically adapted for use with traction drives and to a process for operating traction drives.\n2. Prior Art\nHeretofore, gears and hydraulic devices have been used for power transmission or speed control. A system of drives, called traction drives (rolling friction drive devices), which resort to point contact or line contact between rolling solid bodies has found recognition. These traction drives are particularly useful in industrial machines because they provide high power transmission efficiency without the vibration characteristic of devices using gears and like positive forms of engagement.\nAs the fluid to be interposed between contact parts of such traction drives, there must be selected a lubricant which proves most suitable from the functional point of view. Specifically, the fluid intervening between contact parts is required to possess an ability to undergo reversible glass transition, viscosity increasing, and enhance the efficiency of power transmission between surfaces of rolling contact under high pressure and, on departure from such contact surfaces, immediately return to its original fluid state. It must also be capable of precluding direct contact between metal bodies, preventing metal bodies from seizure, wear and fatigue damage and, similarly to lubricants of all kinds fulfilling important functions of preventing occurrence of rust and elevation of temperature.\nFriction or traction drive devices for transmission of mechanical power have been disclosed in a number of reports in prior technical literature. They are dealt with in full detail in U.S. Pat. Nos. 3,394,603 and 3,411,369; the Journal of Chemical and Engineering Data, Vol. 5, No. 4, p.p. 499-507 (1960), and Hewko et al., in Proceedings of the Symposium on Rolling Contact Phenomena, p.p. 157-185 (1962), Elsevier, Amsterdam, Netherlands, for example.\nAs compositions for use in traction drives, a host of substances have been advanced, including mineral oils (Japanese Patent Publication No. 24,635/1964), mixtures of dialkyl aromatic hydrocarbons with diarylalkanes (Japanese Patent Publication 40,525/1972), polymethylmethacrylate (Japanese Pat. Publication No. 31,828/1973), adamantanes (Japanese Pat. Publications 42,067/1973 and 42,068/1973), polyolefins (Japanese Patents KOKAI 4,766/1971 and 2,229/1972) and alkylnaphthalenes (U.S. Pat. No. 2,549,377), for example. Japanese Patent KOKAI 40,726/1980 has proposed fluids which are obtained by hydrogenating bis-(.alpha.-methylbenzyltoluene) and/or bis-(.alpha.-methylbenzyl)-xylene.\nBesides, proposals abound which concern naphthenic oils having naphthenic rings. These naphthenic oils include dicyclohexylethane (U.S. Pat. No. 3,577,361), dicyclohexylpropane (Japanese Patent Publication 36,105/1978), hydrogenated condensation ring compounds (U.S. Pat. No. 3,411,369), naphthenes containing at least one saturated carbon-containing cyclic ring (U.S. Patent 3,440,894), naphthenes containing at least two saturated carbon-containing cyclic rings (U.S. Pat. Nos. 3,925,217), and mixtures of naphthenes and paraffins (U.S. Pat. Nos. 3,595,796 and 3,595,797), and an oil which is obtained by hydrogenating the product of alkylation of xylene and/or toluene with styrene (Japanese Patent KOKAI No. 43,108/1980), for example.\nU.S. Pat. Nos. 3,440,894 and 3,925,217 mentioned above embrace a wide range of naphthenic compounds and cite a large number of naphthenes by way of example. Most of them are compounds which possess one or two hydrogenated rings. Among the compounds disclosed therein, those which possess three or more hydrogenated rings are limited to a small number, including tercyclohexyls 1,2,3-tricyclohexylpropane and tricyclohexylmethane, for example. Although they have a high traction factor, they have poor practical use because of high viscosity or high crystal-precipitating temperature. It is noted from the prior technical literature that naphthenic oils containing hydrogenated rings are possessed of generally outstanding properties for lubricants, particularly those to be used in traction drives."}
{"text": "1. Field of the Invention.\nThe present invention relates to removable headrests, and more particularly, to such devices that can be readily disassembled for storage.\n2. Description of the Related Art.\nApplicant believes that the closest reference corresponds to U.S. Pat. No. 2,582,571 issued to A. M. Thomas in 1950. The patented invention is not volumetrically efficient and constrain the user to one position. Here, a user can readily transport the disassembled headrest. Also, the supporting boards provide sufficient support without constraining a user's movement. The combination of an elastic member with a non-elastic member provides the necessary gripping support.\nOther patents describing the closest subject matter provide for a number of more or less complicated features that fail to solve the problem in an efficient and economical way. None of these patents suggest the novel features of the present invention."}
{"text": "It is common knowledge that a properly maintained motor vehicle is more dependable, safer, lasts longer, and increases the satisfaction of its owner as compared to a vehicle that is poorly maintained. Furthermore, following a good maintenance schedule helps to preserve the integrity of any warranty that the owner might have on the vehicle. In addition, the residual value of the vehicle is higher at the time of sale or trade-in. Accordingly, systems for implementing and improving the maintenance and repair of vehicles can be extremely beneficial to the owner, family members and others. The manufacturer of a vehicle typically provides a maintenance schedule that includes a recommended timetable or mileage schedule for tire rotation, engine tune-ups, lubrication, oil changes, front end alignments, radiator flushing, and the like. Failure to follow prescribed maintenance can seriously shorten the operative lifetime of the vehicle, and can adversely affect any applicable warranties on the vehicle itself or on the specific components."}
{"text": "The invention relates to signal processing and, more particularly, to a mixed-mode very-large-scale integration (VLSI) architecture and methods for real-time classification of acoustic transients.\nAcoustic transients, short, impulsive bursts of acoustic energy that last between 10 and 100 ms, are a rich source of information in the natural world, and the ability to process them in real time provides a competitive advantage to species. As a result, humans, like other animals, have evolved the ability to quickly and economically process acoustic transients.\nIn the digital world of algorithms and computers, analogous evolutionary forces have caused engineers to develop powerful digital signal processing (DSP) algorithms for classification of acoustic signals on fast DSP engines. Using modern signal processing techniques to recognize acoustic transients in real time is straightforward on modern processors. The challenge of extracting information from signals has been met by powerful mathematical techniques such as wavelet analysis and hidden Markov models. The need for real-time performance has been met by fast and powerful central processing units (CPUs) and special-purpose DSP chips.\nDespite the above, a closer look at the DSP solutions reveals that the burden of real-time processing is borne by increasingly powerful digital processors. The price for success is measured in terms of power dissipation and complexity. Power dissipation scales linearly with the processor\"\"s clock rate. Thus, all else being equal, a 100-MHz processor dissipates 1000 times more power than a 100-kHz processor. Each bump up in clock rate requires increasingly heroic methods to control power dissipation.\nComplexity can be measured by the number of cycles required to perform a calculation and by the surface area of the chip. Increasingly complex algorithms create pressure to increase the complexity of the processor and thus the area of a chip. The problem of greater chip area can be overcome by scaling down the feature size of the fabrication process, but scaling also has physical limits. Moreover, as the feature size scales down, the fabrication process itself becomes more difficult and exacting.\nAll of this contrasts sharply with nature\"\"s solution. The characteristics and advantages of nature\"\"s acoustic processing algorithms are well documented. Natural systems process acoustic information in real time, with precision and reliability, while dissipating minuscule amounts of energy. Nature accomplishes this with slow and unreliable devices, i.e., neurons. Biological hardware has no clock but typical time scales are measured in fractions of milliseconds. In effect, biological hardware runs at a 1- to 10-kHz clock rate.\nIf it were possible to engineer acoustic processors with biological levels of performance and power requirements, a large number of new applications would become feasible. Intelligence based on acoustic pattern recognition could be built into appliances, telephones, and credit cards. Cellular phones could take spoken commands. Smart credit cards could recognize not only passwords, but also the speaker. Digital watches and calculators that run for years on button cells could understand a small vocabulary of spoken words. Self-diagnosing machines could recognize acoustic transients caused by state changes and wear.\nMotivated by the observation that biological systems perform very sophisticated tasks while making low demands on power consumption and component precision, artificial devices can be developed that perform as competently as biological systems while requiring minimal resources. The long-term goal is to build pattern recognition engines whose performance characteristics rival those of biological systems. To be more specific, acoustic processing engines with the following characteristics can be built:\nReal-time operation, so that typical transients are recognized in about 100 ms or less.\nHigh correct classification rates (near 95%) on hundreds of transient classes while achieving low false alarm rates.\nImplementation of such engines with the highly mismatched metal-oxide-silicon (MOS) transistors that are typical in modern analog VLSI fabrication processes (feature size  less than 1.2 xcexcm).\nPower dissipation on the order of a milliwatt or less. This requires subthreshold current-mode circuits. Currents in such circuits are in the 0.1- to 10-nA range, while voltage swings are in the 100-mV range. Clock rates will be tens of kilohertz or less.\nOne solution to the above would be a practical architecture for performing real-time recognition of acoustic transients by means of a correlation-based algorithm. In other words, the algorithm would perform pattern recognition by correlating an incoming signal with a stored template. However, correlation-based algorithms are generally believed to be so computationally intensive that they cannot be used for real-time applications except in conjunction with fast DSP chips.\nTraditionally, correlation in analog VLSI poses two fundamental implementation challenges: first, the problem of template storage; second, the problem of accurate analog multiplication. Both problems can be solved by building sufficiently complex circuits. For example, analog values can be stored by sample-and-hold circuits or by storing digital values and converting them into analog values via digital-to-analog converters. These solutions are generally inferior to digital correlation algorithms because they lead to analog processors that are large compared with their digital counterparts.\nAnother, more compact solution to the template storage problem is to employ the recently developed floating gate devices. Presently, such devices can store precise analog values for years without significant degradation. Moreover, this approach can result in very compact devices. Unfortunately, programming floating gate devices is not particularly easy. It is relatively slow and requires high voltage. Furthermore, the floating gate degrades each time it is reprogrammed. The fabrication of high-quality floating gates also requires advanced fabrication processes that may not be compatible with circuits for other kinds of on-chip processing.\nFinally, even if the analog storage problem could be solved effectively, the problem of building accurate analog-analog multipliers remains. High-quality analog multipliers are notoriously difficult to build. Effective solutions require considerable area on the chip.\nOne solution to the above problems is to sidestep them completely and to develop an algorithm and architecture that require neither analog storage nor analog multiplication. One instance of this approach is to binarize the input and then to correlate it with a binary template. Thus, the correlations can be performed by simple xe2x80x9cXORxe2x80x9d gates. This approach is compact and fast. Thus, there remains a need for analog VLSI devices for real-time classification of acoustic transients that provide a high level of classification and are small and relatively simple to build.\nThe above problems are solved by the invention, a hybrid approach that replaces analog-analog multiplication with analog-binary multiplication. In mixed-mode hardware this operation corresponds to simple binary multiplexing. The algorithm and architecture of the invention are expected to perform a correlation calculation on a special-purpose parallel analog VLSI chip, using a slow clock (about 10 kHz) and with just a few milliwatts of power dissipation.\nMany time-based classification systems compute the correlation of an incoming discrete-time signal or signals with a predetermined template. While for speech and other complex long-term signals it is necessary to perform dynamic time warping (DTW) or similar weighting of the incoming signal, for transients, a simple correlation in the time-frequency domain yields accurate classification results.\nA general form of the simple correlation is                                           c            z                    ⁡                      [            t            ]                          =                              ∑                          m              =              1                        M                    ⁢                      xe2x80x83                    ⁢                                    ∑                              n                =                1                            N                        ⁢                          xe2x80x83                        ⁢                                          x                ⁡                                  [                                                            t                      -                      n                                        ,                    m                                    ]                                            ⁢                                                p                  z                                ⁡                                  [                                      n                    ,                    m                                    ]                                                                                        (        1        )            \nwhere M is the number of frequency channels of the input, N is the maximum number of time bins in the window, x is the array of input signals split into frequency bands, pz is the matrix of template pattern values for pattern z, and t is the current time. This formula produces a running correlation cz(t) of the input array with the template z.\nFor large M and N, this algorithm can be expensive to execute on a DSP in terms of speed and power requirements. However, the approach of the invention lends itself elegantly to low-power parallel analog computation in the form of MOS transistor circuits operating primarily in the subthreshold mode.\nThe algorithm of the invention is a set of modifications to the algorithm of Equation (1). First, the input is normalized, which is essential for the steps which follow. Next, the input and template are transformed into a zero-mean representation by taking the difference of successive samples, the difference between neighboring channels, or some combination thereof. In this form, the template values can be made binary without significantly increasing classification error rates as determined experimentally. Finally, the differencing operation is moved to the output, yielding the simplest possible form of the architecture.\nAssume the input y to the system is a sampled, continuous-valued acoustic signal, split into M frequency bands. The (rectified) energy envelope for each band, denoted x, is computed and then these system inputs are normalize by the function                                           x            ⁡                          [                              t                ,                m                            ]                                =                                    y              ⁡                              [                                  t                  ,                  m                                ]                                                    θ              +                                                ∑                                      k                    =                    1                                    M                                ⁢                                  xe2x80x83                                ⁢                                  y                  ⁡                                      [                                          t                      ,                      k                                        ]                                                                                      ,                            (        2        )            \nwhere xcex8 is a threshold value included to suppress noise during quiet intervals in the input. The normalized input representation is essential to significantly simplify the pattern classifier algorithm for analog inputs, greatly reducing the size and complexity of the hardware implementation but not degrading the classification result.\nSeveral steps are necessary to make the computation less expensive. First, zero-mean transformed input is compared directly to the zero-mean transformed template. This has the effect of subtracting a constant from the result but otherwise does not affect the computation. The template is replaced by the sign of the transformed template. Thus,                                           c            z                    ⁡                      [            t            ]                          =                              ∑                          m              =              1                        M                    ⁢                      xe2x80x83                    ⁢                                    ∑                              n                =                1                            N                        ⁢                          xe2x80x83                        ⁢                                                            x                  xe2x80x2                                ⁡                                  (                                                            t                      -                      n                                        ,                    m                                    )                                            ⁢                                                p                  z                  xe2x80x2                                ⁡                                  [                                      n                    ,                    m                                    ]                                                                                        (        3        )            \nwhere\nxxe2x80x2(txe2x88x92n,m)=x(txe2x88x92n,m)xe2x88x92x(txe2x88x92n,mxe2x88x921) for channel differencing\nand\npxe2x80x2z[n,m]=sign(pz[n,m]xe2x88x92pz[n,mxe2x88x921]).xe2x80x83xe2x80x83(4)\nIt has been shown that binarization of the template has a negligible effect on classification performance. Note that this result does not hold if the input and template are not transformed into a zero-mean representation.\nWith this major simplification of the algorithm, the multiplication can be reduced to a multiplexing function. The required normalization of the input allows a further simplification that does not affect system performance, in which the template values are made binary [0,1] rather than binary [xe2x88x921,1].\nIf time differencing is used for the zero-mean representation, another simplification is possible, noting that the time difference commutes with the summation, so that Equation (1) can be written                                           c            z                    ⁡                      [            t            ]                          =                                            ∑                              m                =                1                            M                        ⁢                          xe2x80x83                        ⁢                                          ∑                                  n                  =                  1                                N                            ⁢                              xe2x80x83                            ⁢                                                x                  ⁡                                      [                                                                  t                        -                        n                                            ,                      m                                        ]                                                  ⁢                                                      p                    z                    xe2x80x2                                    ⁡                                      [                                          n                      ,                      m                                        ]                                                                                -                                    ∑                              m                =                1                            M                        ⁢                          xe2x80x83                        ⁢                                          ∑                                  n                  =                  1                                N                            ⁢                              xe2x80x83                            ⁢                                                x                  ⁡                                      [                                                                                            (                                                      t                            -                            1                                                    )                                                -                        n                                            ,                      m                                        ]                                                  ⁢                                                                            p                      z                      xe2x80x2                                        ⁡                                          [                                              n                        ,                        m                                            ]                                                        .                                                                                        (        5        )            \nIf we let                                                         c              z              xe2x80x2                        ⁡                          [              t              ]                                =                                    ∑                              m                =                1                            M                        ⁢                          xe2x80x83                        ⁢                                          ∑                                  n                  =                  1                                N                            ⁢                              xe2x80x83                            ⁢                                                x                  ⁡                                      [                                                                  t                        -                        n                                            ,                      m                                        ]                                                  ⁢                                                      p                    z                    xe2x80x2                                    ⁡                                      [                                          n                      ,                      m                                        ]                                                                                      ,                            (        6        )            \nthen\ncz[t]=cxe2x80x2z[t]xe2x88x92cxe2x80x2z[txe2x88x921].xe2x80x83xe2x80x83(7)\nIn this forms the time difference can be computed at the output of the correlator rather than at the inputs, yielding three advantages:\n1. Only one difference rather than M differences need to be computed.\n2. Architecturally, the algorithm is less affected by device mismatch when computing the output based on the difference of successive outputs rather than the absolute value of the output.\n3. Most importantly, since the inputs x are rectified, the product xpxe2x80x2 (when pxe2x80x2 is binary [0,1] is always positive and equals either x or zero, which allows us to conveniently implement the entire convolution as an array of simple on/off current switches carrying current in one direction only.\nWith the invention, a high level of classification performance on real-world data can be achieved with no measurable loss of performance in comparison with a traditional, computationally intensive correlation algorithm. Moreover, the mixed-mode architecture of the invention is not significantly harder to implement than the binary-binary correlation. In the acoustic case, the approach requires neither digital-to-analog nor analog-to-digital converters nor the storage of analog values. The algorithm leads to a correlator whose computing surface bears a strong resemblance to conventional dynamic random access memory (RAM)."}
{"text": "Many acoustic modules, such as microphones or speakers, utilize an acoustic membrane to either produce or receive sound. For example, the acoustic membrane of a speaker module may vibrate to produce sound waves that travel into an external environment. However, as the sound waves produced by such an acoustic membrane must be able to travel to the external environment, liquids from the external environment may be able to enter the speaker module and interfere with and/or damage sensitive components. Similarly, the acoustic membrane of a microphone module may need to be exposed to an external environment in order to receive sound waves.\nIn some implementations, various components of such acoustic modules may be made resistant to water and/or other liquids in order to protect sensitive components. However, even when such components are made resistant to liquids, the presence of such liquids may interfere with acoustic operation. For example, the presence of liquid in an acoustic cavity through which acoustic waves must travel either to or from an acoustic membrane may hinder acoustic membrane vibration. Such hindrance may impede proper operation of such an acoustic module even when damage from such liquids is prevented."}
{"text": "1. Field of the Invention\nThe present invention relates to a document conveying apparatus in which an original document is subjected, for a copy operation thereof, to a light irradiation at a predetermined position while the original document is being conveyed through a constant path.\n2. Description of Related Art\nConventionally, in a document conveying apparatus of this kind, when the operator recognizes that a wrong document is inserted or a document sheet is inserted in an inclined state immediately after the insertion of the document, the operator attempts in an ordinary case to draw out the document from the apparatus and to stop the copy operation.\nOn the other hand, in a document conveying apparatus of this kind, a plurality of document detecting means disposed along a document conveying path are so configured to detect a paper jam condition based on a sequential operation thereof according to the conveyance of the document.\nConsequently, when a document is inserted in the apparatus, the paper jam detecting function is automatically started and hence if the document is drawn out, the apparatus regards by itself that a paper jam has occurred. As a result, the operator must take a troublesome action to reset the paper jam detection in each such a situation.\nIt is therefore an object of the present invention to provide a novel document conveying apparatus.\nAnother object of the present invention is to provide a document conveying apparatus in a copying machine in which a document is subjected to a light irradiation at a predetermined position while the document is being conveyed through a constant path."}
{"text": "The invention relates generally to the control of rolling mills in a manner that offsets the cylic effects of roll assemblies on the thickness of material being rolled, hereinafter referred to as eccentricity.\nA system for controlling the effects of roll eccentricity in a rolling mill is disclosed in U.S. Pat. No. 4,222,254 to King et al, the disclosure of this patent being incorporated herein by reference. An improvement in the King et al system is disclosed in pending U.S. Ser. No. 591,277, filed Mar. 19, 1984 in the name of Mark Puda, now U.S. Pat. No. 4,531,392. The disclosure of the Puda patent is also incorporated herein by reference.\nBoth the King et al and Puda disclosures employ a process by which the eccentricity of one or more roll assemblies of a mill is inferred by making mathematical estimates of eccentricity. Estimates of eccentricity are made because the actual eccentricity of the rolls is not readily observable.\nThe estimates made in the King et al patent and the Puda patent are based on measurements of changes in the force at which the rolls of the mill engage the material being rolled in the mill, measuring changes in the stretch or compression of the mill housing during the rolling process and measuring changes in the position of the actuator mechanisms (mechanical screws or hydraulic cylinders) that control the working space and rolling force of the mill. These measurements are then used in a gaugemeter equation to estimate changes in the thickness of the material exiting the mill due to roll eccentricity.\nThe estimating process is continuous until the cyclic component in exit thickness due to roll eccentricity is reduced to a minimum or zero amount and the estimate of eccentricity reaches a steady state value representing the true eccentricity of the rolls.\nAs discussed in the Puda patent, the load cells and actuator mechanisms of the mill do not respond as fast as the occurrence of the eccentricity disturbance at high travel speeds of the material through the mill. Hence, the resulting control signals of the system can contain a delay in their response to eccentricity. The Puda disclosure cares for this by a process of developing phase compensated gains that are used to modify the amplitude and phase of the estimate of eccentricity calculated by the means disclosed in the King et al patent. The control signals directed to the actuator mechanisms of the mill are thereby phase corrected to care for the delay in the response of the load cells and actuator mechanism."}
{"text": "Storing and safeguarding electronic content is of paramount importance in modern business. Accordingly, various systems may be employed to protect such electronic content.\nThe use of solid-state storage devices is increasing in popularity. A solid state storage device is a content storage device that uses solid-state memory to store persistent content. A solid-state storage device may emulate (and therefore replace) a conventional hard disk drive. Additionally/alternatively, a solid state storage device may be used within a cache memory system. With no moving parts, a solid-state storage device largely eliminates (or greatly reduces) seek time, latency and other electromechanical delays and failures associated with a conventional hard disk drive."}
{"text": "This invention relates to an electrochemical measuring method utilizing photochemical reactions and an apparatus therefor and more particularly relates to a measuring method suitable for detection of components in solutions and analysis of mechanism of photochemical reactions.\nHitherto, as methods for detecting substances contained in a solution at a high sensitivity by irradiation of light to the solution, there have been known such methods as making the detection by measuring absorption of the irradiated light per se, by measuring fluorescence emitted due to excitation with the irradiated light, by using the thermal lens effect caused by thermal expansion due to light absorption heat and by measuring light-acoustic effect, etc. These methods have the problems that they are liable to be affected by noise caused by external stray light, radiant heat or other noise because the detection is carried out with light, heat or sound."}
{"text": "Research is currently being done to develop new sensors (such as environmental sensors) which can be produced at lower cost with improved sensitivity when compared to existing sensors.\nOne or more aspects/embodiments of the present disclosure may or may not address one or more of these issues.\nThe listing or discussion of a prior-published document or any background in this specification should not necessarily be taken as an acknowledgement that the document or background is part of the state of the art or is common general knowledge."}
{"text": "Flow control valves are used for controlling fluid flows in a variety of applications. For instance, in vapour compression systems, such as refrigeration systems, air condition systems or heat pumps, flow control valves are, e.g., used for controlling a supply of refrigerant to an evaporator. In this case the flow control valve may be in the form of an expansion valve. Flow control valves may also be used for controlling a flow of refrigerant in a suction line of a vapour compression system, i.e. downstream relatively to an evaporator. In this case the flow control valve functions as a suction valve.\nU.S. Pat. No. 7,240,694 B2 discloses a flow control valve for a refrigeration system. The flow control valve comprises a valve body defining a fluid inlet section, a fluid outlet section, and a connecting port interconnecting the inlet and outlet sections. A valve unit is movable to selectively opening and closing the valve port, thereby controlling the fluid flow between the fluid inlet section and the fluid outlet section.\nThe flow control valve of U.S. Pat. No. 7,240,694 B2 is made from a number of separate parts. When the flow control valve is assembled these separate parts are attached to the valve, one by one, at the operating site of the flow control valve. This is cumbersome and makes it difficult to manage the individual parts of the flow control valve."}
{"text": "Throughout this application, various publications are referenced in parentheses. Citations for these references may be found at the end of the specification immediately preceding the claims. The disclosures of these publications in their entireties are hereby incorporated by reference into this application to more fully describe the state of the art to which this invention pertains.\nSingle nucleotide polymorphisms (SNPs), the most common genetic variations in the human genome, are important markers for identifying disease genes and for pharmacogenetic studies (1, 2). SNPs appear in the human genome with an average density of once every 1000-base pairs (3). To perform large-scale SNP genotyping, a rapid, precise and cost-effective method is required. Matrix-assisted laser desorption/ionization time-of-flight mass spectrometry (MALDI-TOF MS) (4) allows rapid and accurate sample measurements (5–7) and has been used in a variety of SNP detection methods including hybridization (8–10), invasive cleavage (11, 12) and single base extension (SBE) (5, 13–17). SBE is widely used for multiplex SNP analysis. In this method, primers designed to anneal immediately adjacent to a polymorphic site are extended by a single dideoxynucleotide that is complementary to the nucleotide at the variable site. By measuring the mass of the resulting extension product, a particular SNP can be identified. Current SBE methods to perform multiplex SNP analysis using MS require unambiguous simultaneous detection of a library of primers and their extension products. However, limitations in resolution and sensitivity of MALDI-TOF MS for longer DNA molecules make it difficult to simultaneously measure DNA fragments over a large mass range (6). The requirement to measure both primers and their extension products in this range limits the scope of multiplexing.\nA high fidelity DNA sequencing method has been developed which uses solid phase capturable biotinylated dideoxynucleotides (biotin-ddNTPs) by detection with fluorescence (18) or mass spectrometry (19), eliminating false terminations and excess primers. Combinatorial fluorescence energy transfer tags and biotin-ddNTPs have also been used to detect SNPs (20).\nFalse stops or terminations occur when a deoxynucleotide rather than a dideoxynucleotide terminates a sequencing fragment. It has been shown that false stops and primers which have dimerized can produce peaks in the mass spectra that can mask the actual results preventing accurate base identification (21).\nThe present application discloses an approach using solid phase capturable biotin-ddNTPs in SBE for multiplex genotyping by MALDI-TOF MS. In this method primers that have different molecular weights and that are specific to the polymorphic sites in the DNA template are extended with biotin-ddNTPs by DNA polymerase to generate 3′-biotinylated DNA extension products. The 3′-biotinylated DNAs are then captured by streptavidin-coated magnetic beads, while the unextended primers and other components in the reaction are washed away. The pure DNA extension products are subsequently released from the magnetic beads, for example by denaturing the biotin-streptavidin interaction with formamide, and analyzed with MALDI-TOF MS. The nucleotide at the polymorphic site is identified by analyzing the mass difference between the primer extension product and an internal mass standard added to the purified DNA products. Since the primer extension products are isolated prior to MS analysis, the resulting mass spectrum is free of non-extended primer peaks and their associated dimers, which increases the accuracy and scope of multiplexing in SNP analysis. The solid phase purification system also facilitates desalting of the captured oligonucleotides. Desalting is critical in sample preparation for MALDI-TOF MS measurement since alkaline and alkaline earth salts can form adducts with DNA fragments that interfere with accurate peak detection (21)."}
{"text": "The management and treatment of myocardial infarction has changed dramatically since the first half of the 20th century, progressing from an era of bed rest and observation, to an emphasis on technology, including hemodynamic monitoring and balloon catheters, to an increased focus on thrombolytic therapy. (Antman and Braunwald, “Acute Myocardial Infarction” in Heart Disease, A Textbook of Cardiovascular Medicine, 61 edition, vol. 2, Braunwald et al., eds, 2001, W.B. Saunders Company, Philadelphia). Therapeutic approaches to treating cardiovascular diseases have evolved tremendously in the last 100 years accompanied by greater understanding of the underlying pathology.\nAlmost all myocardial infarctions result from coronary atherosclerosis, generally with superimposed coronary thrombosis. Slowly accumulating plaques can be asymptomatic due to the development of collateral vessels. However, atherosclerotic plaques, especially those rich in lipids, are prone to abrupt plaque rupture. Plaque rupture and associated endothelial injury cause the release of mediators such as thromboxane A2, serotonin, adenosine diphosphate, thrombin, platelet activating factor, tissue factor and oxygen-derived free radicals. These mediators promote platelet aggregation and mechanical obstruction often leading to thrombus formation which interferes with blood flow and oxygen supply. Persistent and severe interferences with myocardial oxygen supply can lead to acute myocardial infarction. (See, Rioufol et al., 2002, Circulation 106:804, Timmis, 2003, Heart 89:1268-72).\nThe mainstay of atherosclerotic pharmacotherapy has been chronic therapy to prevent or slow the development of atherosclerotic plaques primarily by focusing on lowering LDL or “bad cholesterol” as a therapeutic endpoint. Statin therapy, for example, has greatly contributed to improved cardiovascular health; however, adverse effects such as rhabdomyolysis, remain an impediment. Furthermore, statins do little in an acute situation, for example, to reduce vulnerable, unstable atherosclerotic plaque during an ischemic episode. Acute treatment has largely relied on thrombolytics (such as tPA) and surgical intervention such a percutaneous transluminal coronary angioplasty (PTCA) and coronary artery bypass graft (CABG). While thrombolytics provide relief by decreasing or eliminating an occluding thrombus, they do not alter the underlying pathology. Interventions such as PTCA carry their own risks and are often unsuitable for patients in acute conditions. Hence current pharmacologic therapies do little to help patients once unstable plaque presents as a risk. (See, Newton and Krause 2002, Atherosclerosis S3:31-38).\nYet, despite the improved understanding of the pathophysiology of myocardial infarction and developments in atherosclerotic pharmacotherapy, safe and effective treatment modalities which have a fast onset of action to allow for treatment in the acute phase, and which do not have serious side effects when used for long term treatment, are still desired."}
{"text": "1. Field of the Invention\nGenerally, the present disclosure relates to the formation of integrated circuits, and, more particularly, to the formation of a dielectric interlayer between and over circuit elements including closely spaced lines, such as gate electrodes, polysilicon interconnect lines and the like.\n2. Description of the Related Art\nDuring the fabrication of integrated circuits, a large number of circuit elements are formed on a given chip area according to a specified circuit layout. Generally, a plurality of process technologies are currently practiced, wherein, for complex circuitry, such as micro-processors, storage chips and the like, MOS technology based on silicon is currently the most promising approach due to the superior characteristics in view of operating speed and/or power consumption and/or cost effectiveness. During the fabrication of complex integrated circuits using MOS technology, millions of transistors, i.e., N-channel transistors and/or P-channel transistors, are formed on a substrate including a crystalline semiconductor layer, such as a silicon-based layer. A MOS transistor, irrespective of whether an N-channel transistor or a P-channel transistor is considered, comprises so-called PN junctions that are formed by an interface of highly doped drain and source regions with a lightly doped channel region disposed between the drain region and the source region. The conductivity of the channel region, i.e., the drive current capability of the conductive channel, is controlled by a gate electrode that comprises a line-like portion and is formed above the channel region and separated therefrom by a thin insulating layer.\nTypically, the circuit elements, such as the MOS transistors, capacitors, resistors and the like, are formed in a common layer, which will be referred to hereinafter as a device layer, whereas the “wiring,” i.e., the electrical connection of circuit elements according to the circuit design, may be accomplished only to a certain degree by means of polysilicon lines and the like within the device layer so that one or more additional “wiring” layers formed over the device layer may be required. These wiring layers include metal lines embedded into an appropriate dielectric material, such as silicon dioxide, silicon nitride and the like, or, in advanced devices, low-k materials having a permittivity of 3.0 or less are used. The metal lines and the surrounding dielectric material will be referred to hereinafter as a metallization layer. Between two stacked adjacent metallization layers and also between the device layer and the first metallization layer, respective dielectric interlayers are formed through which metal-filled openings are formed to establish the electrical connection between metal lines or between circuit elements and metal lines. In typical applications, the dielectric interlayer separating the device layer from the first metallization layer is essentially formed from silicon dioxide that is deposited above a dielectric etch stop layer by well-established plasma enhanced chemical vapor deposition (PECVD) techniques, which enable the formation of a smooth and dense silicon dioxide film with sufficient conformality at moderately high deposition rates. Due to the continuous device scaling resulting in gate lengths of MOS transistors on the order of 50 nm or less, the distances between neighboring circuit elements, such as polysilicon lines, gate electrodes and the like, are also reduced and have now reached in modern CPUs approximately 200 nm and less, which translates into approximately 100 nm or less for the space width between the densely packed polysilicon lines. It turns out, however, that the gap-fill capabilities of well-established high rate PECVD techniques for the deposition of silicon nitride, which is frequently used as material for the etch stop layer, and silicon dioxide, which is often used as interlayer dielectric, may no longer suffice to reliably form a dielectric interlayer, thereby requiring a fill technique providing enhanced fill capabilities as will be described in more detail with reference to FIGS. 1a-1b. \nIn FIG. 1a, a semiconductor device 100 comprises a substrate 101 that may be a bulk silicon substrate or a silicon-on-insulator (SOI) substrate having formed thereon a device layer 102 including, for instance, a silicon-based layer 110, in and on which is formed a structure 103 that may comprise closely spaced polysilicon lines 104. The device layer 102 may represent a substantially crystalline silicon region in which and on which circuit elements, such as field effect transistors, capacitors and the like, are formed. The structure 103 may represent an area having a plurality of dense polysilicon lines, or the lines 104 may represent portions of gate electrodes of transistor elements. The lines 104 may have formed on sidewalls thereof corresponding spacer structures 105, as are typically used for forming gate electrode structures. The spacer structures 105 may include a plurality of spacers, such as an offset spacer 105A and one ore more “outer” spacers 105C, and a liner 105B that may act as an etch stop layer during an etch process for forming the respective spacers 105C. The structure 103 further comprises an etch stop layer 109, typically comprised of silicon nitride, that is formed over the device layer 102 to cover the layer 110 and the line structure 103. A silicon dioxide layer 107 is formed above the etch stop layer 109 so as to completely enclose the line structure 103.\nA typical conventional process flow for forming the device 100 as shown in FIG. 1a may include the following processes. After fabrication processes to form circuit elements, such as transistors, capacitors and the line structure 103, which include well-established lithography, deposition, etch, implantation and other techniques, the etch stop layer 109 is formed, typically by PECVD, since PECVD of silicon nitride may be accomplished at moderately low temperatures of less than approximately 600° C., which is compatible with preceding manufacturing processes and materials, such as metal silicides and the like. In many conventional techniques, the etch stop layer 109 may be provided with a high intrinsic stress level so as to act as a strain-inducing source for creating a strain in an area 108 located below the lines 104. When the lines 104 represent gate electrodes, the area 108 may be considered as a channel region of a transistor, in which the induced strain may result in a modified charge carrier mobility. For example, for a standard crystallographic orientation of the semiconductor layer 110, that is, when the layer 110 represents a silicon-based material having a surface orientation (100) with the channel length oriented along a <110> direction, a compressive strain in the area 108 may result in an improvement of the hole mobility while a tensile strain may result in an improvement of the electron mobility. The enhanced charge carrier mobility thus directly translates into enhanced transistor performance with respect to current drive capability and operating speed. In order to selectively enhance the transistor performance, the etch stop layer 109 may be deposited on the basis of appropriately selected process parameters so as to obtain the desired degree and type of intrinsic stress. For example, silicon nitride may be deposited by PECVD with high tensile or compressive stress, depending on the deposition parameters. Moreover, well-established process sequences may be used to selectively form portions of the etch stop layer 109 with a different type of intrinsic stress above different transistors in order to enhance the performance of both N-type transistors and P-type transistors.\nAs previously discussed, the ongoing shrinkage of feature sizes also means that a distance between neighboring circuit elements, such as a distance 111 between the closely spaced lines 104, is reduced and may be as low as approximately 100 nm, or the distance 111 may even be as small as 30 nm and even less for CPUs of the 90 nm technology node. Hence, any deposition techniques for forming a dielectric layer for embedding the line structure 103 with open spaces therebetween have to meet the requirements of an appropriate fill capability so as to reliably and completely fill the empty spaces between the densely spaced lines 104. By means of PECVD process recipes for silicon nitride, the layer 109 may be deposited in a more or less conformal fashion with a thickness in the range of approximately 10-100 nm, wherein possibly different types of intrinsic stress may be provided above respective portions of the structure, thereby requiring sophisticated deposition and patterning strategies, in particular when the creation of voids 106a is to be suppressed.\nThereafter, the silicon dioxide layer 107 is deposited, which in less critical applications is typically done by PECVD on the basis of precursors TEOS (tetra-ethyl-ortho-silicate) and oxygen, since PECVD, contrary to thermal TEOS chemical vapor deposition (CVD), allows the deposition of silicon dioxide in a moderately conformal manner—yet with significantly less gap filling qualities compared to thermal CVD—with relatively high mechanical stability at temperatures below 600° C. at high deposition rates, which provides a high production yield.\nHowever, with the distance 111 approaching approximately 30 nm and even less, it turns out that the fill capabilities of well-established PECVD techniques for depositing silicon dioxide having superior material characteristics on the basis of TEOS and oxygen may not be adequate to completely fill the empty spaces between the lines 104, thereby possibly creating voids 106b, which may lead to severe reliability concerns during the further processing of the semiconductor device 100, i.e., during the fabrication of contacts providing an electrical connection between individual elements of the structure 103 to a metallization level to be formed. Moreover, it should be noted that the silicon dioxide layer 107 has a certain topography caused by the underlying structure of the device layer 102, for instance, by the line structure 103, which may jeopardize subsequent manufacturing processes, such as a photolithography step for forming contact openings to underlying portions of circuit elements located in the layer 110 or on the lines 104. Consequently, a standard process flow requires that the silicon dioxide layer 107 be planarized, typically by chemical mechanical polishing (CMP), wherein excess material of the silicon dioxide layer 107 is removed by chemical and mechanical interaction with a slurry and a polishing pad so as to finally obtain a substantially planarized surface of the silicon dioxide layer 107. The CMP process itself is a highly complex process and requires sophisticated process recipes, which significantly depend on the characteristics of the silicon dioxide layer 107, such as density, mechanical stress, water contents and the like. Hence, a great deal of effort is required to develop corresponding process recipes for reliable and reproducible CMP processes for PECVD TEOS silicon dioxide, as this material is frequently used for a dielectric interlayer in silicon-based semiconductor devices and even in devices formed from other semiconductors.\nFor this reason, the dielectric layer 107 formed on the silicon nitride layer 109 may be deposited by a different deposition technique having a significantly enhanced gap filling capability to avoid the creation of the voids 106b. Hence, the silicon dioxide layer 107 may be formed by a thermal CVD process on the basis of TEOS and ozone, which generates a silicon dioxide film exhibiting excellent gap filling capabilities, that is, this deposition technique provides even a “flow”-like behavior, thereby reliably filling the empty spaces between the lines 104. In view of the film and deposition characteristics, the thermal CVD process is typically performed at significantly higher pressures compared to the plasma enhanced deposition technique, for example, in the range of 200-760 Torr, and is therefore denoted as sub-atmospheric chemical vapor deposition (SACVD). However, the material and process characteristics of the SACVD oxide may differ significantly from the PECVD oxide, as, for instance, the layer 107 formed by SACVD may tend to incorporate moisture more readily and also exhibit an increased rate of out-gassing compared to PECVD oxide. Furthermore, the deposition rate is lower, resulting in a reduced throughput. For these reasons, the layer 107 is provided as an intermediate material used as a gap fill material and thereafter a further silicon dioxide layer 107A may be deposited by PECVD to provide the desired deposition rate and enhanced material characteristics for at least the upper portion of the interlayer dielectric material. Thus, during the further processing, for instance the planarization of the interlayer dielectric material 107A, well-established process techniques may be used, while, however, the inferior material characteristics of the SACVD oxide may have an inverse effect on the overall reliability of the final interlayer dielectric material and thus on the structure 103.\nFIG. 1b schematically illustrates the semiconductor device 100 according to another illustrative example in which the deposition process having the desired high gap filling capability may result in a high degree of non-uniformity during the further processing of the device 100. As shown, the device 100 may comprise the etch stop layer in the form of a first portion 109A having a high intrinsic stress level, for instance, a high compressive stress, while a second portion 109B may have a high intrinsic stress level of opposite behavior, such as a tensile stress. As previously explained, the lines 104 of the structure 103 may represent gate electrode structures of transistors, in which an appropriately selected type of strain in the respective channel regions 108 may provide enhanced transistor performance, as previously explained. When forming the portions 109A, 109B, respective deposition parameters may be adjusted, such as the deposition pressure, temperature, precursor flow rate, ion bombardment and the like, in order to obtain the desired high intrinsic stress levels. For example, according to well-established process recipes, a stressed dielectric material may be deposited in a highly conformal manner and a portion thereof may then be removed to obtain, for instance, the portion 109A. Thereafter, the dielectric material may be deposited with the opposite intrinsic stress level to that of the portion 109B while an unwanted part thereof may be removed from above the portion 109A, thereby obtaining the configuration as shown in FIG. 1b. \nDuring these manufacturing processes, the respective deposition parameters may also be selected so as to obtain a highly conformal deposition behavior in order to substantially avoid the creation of any voids between the densely spaced lines 104. Thereafter, the interlayer dielectric material 107 or a portion thereof may be deposited on the basis of the sub-atmospheric deposition process, as previously described, in order to ensure a reliable filling of the spaces between the lines 104. It turns out, however, that the growth rate during this deposition process may be different for a material having a high compressive stress and a tensile-stressed dielectric material, thereby resulting in a different layer thickness of the interlayer dielectric material 107 above the portions 109A, 109B. Consequently, during the further processing, for instance, when providing a further interlayer dielectric material, such as the material 107A, planarizing the resulting surface topography and the like, an increased degree of process non-uniformity may be encountered, which may also result in respective device non-uniformities, for instance, in view of a reduced planarity and the like.\nThus, although the enhanced gap filling capabilities of the sub-atmospheric deposition technique for silicon dioxide may be highly advantageous with respect to avoiding structure irregularities, in particular in densely packed line structures and gate electrodes, the inferior material characteristics, possibly in combination with deposition specific non-uniformities, may result in a reduced reliability and increased device irregularities, in particular for highly scaled semiconductor devices.\nThe present disclosure is directed to various techniques and devices that may avoid, or at least reduce, the effects of one or more of the problems identified above."}
{"text": "Electromechanical devices generally comprise a class of devices that combine electrical and mechanical parts. There are many types of electromechanical devices, and examples include microelectromechanical (MEM) devices, microelectromechanical systems (MEMS), microsystems (MST), nanoelectromechanical systems (NEMS), sensors, transducers, actuators and switches. Electromechanical devices having planar configurations offer several advantages over nonplanar configurations, including reduced size, lower power consumption, and lower fabrication costs.\nThe two most widely used techniques for fabricating planar electromechanical devices are surface micromachining (SM) and bulk micromachining (BM). While SM defines a structure by deposition and etching of different structural layers, BM defines a structure by selectively etching inside a substrate. The differences in these two manufacturing processes results in differences in structures and properties of devices fabricated thereby. For example, due to the conformal nature of SM, which involves successive depositions of metals and dielectrics, nonplanar structures also known as step beams are formed. Switches embodying these step beams are susceptible to latching or friction when a switch's cantilever conforms to its underlying electrical contact. In contrast, BM, which can include wafer bonding, yields planar structures.\nFurther, BM uses single crystal materials, which are superior to the deposited films used in SM. For example, single crystal substrates tend to have fewer crystal lattice defects than thin films. In addition, the mechanical properties of single crystal substrates (e.g., Young's modulus and Poisson's ratio) are highly repeatable, which again facilitates fewer crystal lattice defects. In contrast, the mechanical properties of thin films vary widely with the conditions under which such films are processed. Furthermore, while single crystal substrates are substantially free of built-in stresses, deposited thin films may include a variety of built-in compressive and tensile stresses that detrimentally affect manufacturing and performance. Due to these shortcomings, surface micromachined switches may develop stress concentration points during switch actuation which, over time, can lead to device failure. Similarly, contact dimples formed on switches using SM technology are prone to failure due to delaminations occurring between the thin film layers during extended periods of switch actuation.\nIn BM processing technology, the most popular substrate is silicon wafers due to the favorable anisotropic properties of silicon in which its crystal structure is arranged in lines and planes. Because of this structural arrangement, etching can be selectively performed on specific lines and planes that have relatively weak bonds. However, given the inferior insulation properties of silicon vis-a.-vis other materials, RF planar switches comprising silicon exhibit relatively low isolation and thus high insertion losses.\nRF switches are widely used in a variety of applications including, for example, telecommunication applications. In this regard, RF switches are extremely important building blocks for reconfigurable RF communication systems. In one application, the use of planar RF switches can reduce the overall size, weight and cost of switch matrices on satellites. In other applications, planar RF switches can be incorporated into software programmable radio systems, reconfigurable antennas for radar, and antennas for mobile communications.\nAs can be seen, there exists a need in the art for improved methods and apparatus for planar RF switch technology offering a durable switch made from a single crystal in which the switch has high isolation, low insertion losses and highly repeatable mechanical properties. The embodiments of the present disclosure answer these and other needs."}
{"text": "Machine virtualization is commonly used in various computing environments, such as in data centers and cloud computing. A typical virtualized computing system comprises multiple compute nodes that are interconnected by a high-speed network. The compute nodes run Virtual Machines (VMs) that consume physical resources such as Central Processing Unit (CPU), memory, persistent storage and networking resources. Some computing systems use distributed File Systems (FSs) for persistent storage of files. Examples of distributed FSs include Google File System (GFS), Hadoop distributed file system (HDFS), Ceph, General Parallel File System (GPFS) and Lustre.\nU.S. Pat. Nos. 8,863,124 and 8,601,473, whose disclosures are incorporated herein by reference, describe techniques for implementing I/O and storage device management in a virtualization environment. A Service VM is employed to control and manage various types of storage devices, including directly attached storage in addition to networked and cloud storage."}
{"text": "1. Field of the Invention\nThe invention generally relates to arrangements for suppressing or rejecting “images” (energy at harmonic frequencies) in communications systems. More particularly, the invention relates to arrangements that use a current switching technique to combine a digital-to-analog converter (DAC) function with a filtering function to significantly reduce harmonic images.\n2. Related Art\nFIG. 1 shows a conventional transmission arrangement, in which digital data are input to a digital-to-analog converter (DAC) 100. DAC 100 outputs an analog signal that changes value at discrete times determined by a clock signal CK as it clocks D-type flip-flops (not specifically shown) at the DAC's input. Being the result of discrete time clocking, the DAC's analog output signal has sharp vertical edges in the time domain.\nUndesirably, the presence of sharp edges in the DAC's analog output signal implies there is significant energy at harmonic frequencies (see FIG. 5A; which assumes a CK frequency of 1 GHz). To mask interfering emissions such as those at harmonic frequencies (as effectively mandated, for example, by the U.S. Federal Communications Commission, FCC), a low pass reconstruction filter 102 has conventionally been employed. Such a filter lowers the higher-frequency content of signals before sending them to a modulator 104 for generating a radio frequency (RF) modulated signal.\nThe reconstruction filter is usually an analog filter that consumes a substantial amount of power, has large process/temperature/voltage variations (approximate ±30% for integrated passive R/C components), and is very difficult to design when the DAC input/output signal frequency is high. Furthermore, the closeness of the harmonic frequencies to the desire pass band, and the prevalence of higher energy content at lower harmonics (see FIG. 5A) have required filter 102 be complex.\nConventionally, higher order analog filters (such as Butterworth filters) have been used to implement filter 102. Undesirably, such complex filters are extremely difficult to design when the corner filter frequency is high. Also, operational amplifiers required in analog filters consume large amounts of power. Moreover, the unity gain bandwidth is often impossible to achieve with a high filter corner frequency. The accuracy of the corner frequency adjustment may sometimes be achieved by trimming or by calibration, but either of these approaches significantly increases design complexity.\nAccordingly, there is a need in the art for an arrangement that simply and efficiently converts digital signals into analog signals that are lower in harmonic content, thereby avoiding or minimizing the need for the costly filters found in conventional systems."}
{"text": "1. Field of the Invention\nThe present invention relates to an item display method and apparatus. More particularly, the present invention relates to a method and apparatus that display items according to a user gesture.\n2. Description of the Related Art\nFor playback of pre-stored digital content such as songs, images, games, broadcasts or moving images, an item display apparatus typically displays a list of items on the screen representing digital content. Such an item display apparatus may provide various input schemes adapted to enable the user to easily locate and select a desired item.\nWhen a large amount of content is present, the item display apparatus may divide and even sub-divide many items into several categories for easier item management. The item display apparatus may provide a tree-type user interface, through which the user may traverse categories of items. Using the tree-type user interface, the user may search the categories for an item to be executed in a hierarchical manner. However, in devices such as mobile terminals having a small screen, as the number of categories and the layers thereof increase, users may have to frequently make screen transitions to access a desired item."}
{"text": "Today's personal computers (PCs) and workstations are experiencing a rapid growth of accelerated clocking and computational speeds. PC clock speeds have progressed from the Intel Corporation 486™ microprocessor speeds of 60 and 90 MHz, to the present Pentium III™ clock speeds in excess of 600 MHz. Moore's Law generally predicts a doubling of computing power and circuit complexity every year and a half or so.\nThe downside corollary to Moore's Law, however, is that with doubling the number of devices in an integrated circuit (IC) consequently raises the amount of heat generated requiring dissipation. As an integrated circuit drives current between transistors it consumes power, producing waste heat that eventually transfers outward through the chip from the surface of the die. Generally, a PC chip designed for commercial use can withstand up to 150° C. Exceeding that limit, however, will cause the chip to make errors in its calculations, or perhaps fail completely.\nCurrent solutions in heat dissipation for chips include heat spreaders, heat sinks and fans. Heat spreaders, which generally are made of a tungsten-copper alloy and are placed directly over a chip, have the effect of increasing the chip's surface area, allowing more heat to be vented upward. Similarly, heat sinks spread the heat upward through fins or folds, which are vertical ridges or columns that allow heat to be conducted in three dimensions—length, width, and height, as opposed to the two-dimensional length and width of heat spreaders.\nFans within a computer housing can further aid heat dissipation from the chip or heat sink surface by convection. The amount of heat a fan dissipates away from a chip depends on the volume of air the fan moves, the ambient temperature, and the difference between the chip temperature and the ambient temperature.\nThe miniaturizing of integrated circuits have generally allowed for a reduction of operating voltages, resulting in lower heat production. However, chip shrinkage also means that heat-generating devices are packed closer together. Thus, the “power density” or the amount of heat concentrated at particular spots across the chip may begin to climb. As a consequence, heat is generated faster than it can be dissipated as higher clock speeds are demanded.\nFor example, PCs with a 486™ microprocessor drew generally 12 to 15 watts, primarily concentrated in the processor. A power supply with an embedded fan was typically sufficient to circulate air and cool the inside of the PC chassis, while a passive heat sink could cool the processor. On the other hand, Pentium™ processors from Intel Corporation consumed about 25 watts, thus requiring more cooling means than the passive heat sink for the processor alone. Similarly, the Pentium II™ processor consumes about 40 watts, while future processors like the 64-bit Merced™ may consume up to 65 watts. Other transistor-laden components must also contend with increasing heat generation: add-on cards; chipsets; graphics chips; and high-performance dynamic random access memory (DRAM).\nThere is, consequently, a need for improved heat dissipation from integrated circuits."}
{"text": "Radio, television, satellite, and other broadcasters often devote a significant amount of time and effort to scheduling media content and spots for broadcast. In some circumstances, the same content is delivered for broadcast over two different media channels, for example when a radio station broadcasts media content over the air, and streams the same media content via an Internet radio station. In some cases, even though both stations are broadcasting the same programming content, commercials, or spots, broadcast during particular spot blocks, may vary between stations. Broadcasting different spots during a common spot block can greatly increase the complexity involved in scheduling spots. It can also be difficult to make sure that the length of the spot blocks on both of the stations remain consistent, and that the spot blocks remain synchronized. Conventional technologies do not always provide the most effective way to handle the scheduling and synchronization issues that can arise in various broadcasting situations."}
{"text": "Global communications networks such as the Internet are now ubiquitous with an increasingly larger number of private and corporate users dependent on such networks for communications and data transfer operations. As communications security improves, more data can be expected to traverse the global communications data backbone between sources and destinations (typically, server hosts) placing increasing demands on those entities that handle and store data. Such increased demands are typically addressed at the destination by adding more switching devices and servers to handle the load. However, this can be an expensive proposition in terms of hardware, software, setup, and administration.\nNetwork load-balancers provide client access to services hosted by a collection of servers (herein known as “hosts”). Clients connect to a load-balancer, which transparently (to the clients) forwards them to a host according to a set of rules. This general load balancing context includes the following: packets form sequences, called sessions; sessions should be allocated among the available hosts in a “balanced” manner; and, every packet of each session should always be directed to the same host, so long as the host is alive (the latter is known as “session affinity”).\nThis problem is most often solved through the use of a single monolithic load-balancer that monitors the status (liveness/load) of the hosts and maintains state in the form of a table of all active sessions. When a new session arrives, the load-balancer selects the least-loaded host that is available and assigns the session to that host. In order to provide session affinity, the load-balancer must “remember” this assignment (routing) decision by adding an entry to its session table. When subsequent packets for this session arrive at the load-balancer, a single table lookup determines the correct host. However, an individual load-balancer can be both a single point of failure and a bottleneck; the size of its session table (and thereby the amount of state maintained) grows with increased throughput and routing decisions for existing session traffic require a state lookup (one per packet). Circumventing these limitations requires multiple load-balancers working in tandem (scale-out), and/or larger, more powerful load-balancers (scale-up). However, scaling-out these load balancing devices is complicated, due most notably to the need of maintaining consistent state among the load-balancers. Likewise, scaling them up is expensive, since cost versus throughput in fixed hardware is non-linear (e.g., a load-balancer capable of twice the throughput costs significantly more than twice the price)."}
{"text": "EP 238 922 B1, which is incorporated by reference, discloses a position sensor that operates on the principle of measuring linear displacement on the basis of a permanent magnetic linear contactless displacement, known as PLCD. Such a position sensor is also referred to as a linear inductive position sensor, known as a LIPS."}
{"text": "One of the most important and fundamental processes in biology is the division of cells mediated by the cell cycle. This process ensures the controlled production of subsequent generations of cells with defined biological function. It is a highly regulated phenomenon and responds to a diverse set of cellular signals both within the cell and from external sources. A complex network of tumor promoting and suppressing gene products are key components of this cellular signaling process. Over expression of the tumor promoting components or the subsequent loss of the tumor suppressing products will lead to unregulated cellular proliferation and the generation of tumors (Pardee, Science 246:603-608, 1989).\nCyclin dependent kinases (cdks) play a key role in regulating the cell cycle machinery. These complexes consist of two components: a catalytic subunit (the kinase) and a regulatory subunit (the cyclin). To date, six kinase subunits (cdk 1-7) have been identified along with several regulatory subunits (cyclins A-H). Each kinase associates with a specific regulatory partner and together make up the active catalytic moiety. Each transition of the cell cycle is regulated by a particular cdk complex: G1/S by cdk2/cyclin E, cdk4/cyclin D1 and cdk6/cyclinD2; S/G2 by cdk2/cyclin A and cdk1/cyclin A; G2/M by cdk1/B. The coordinated activity of these kinases guides the individual cells through the replication process and ensures the vitality of each subsequent generation (Sherr, Cell 73:1059-1065, 1993; Draetta, Trends Biochem. Sci. 15:378-382, 1990)\nAn increasing body of evidence has shown a link between tumor development and cdk related malfunctions. Over expression of the cyclin regulatory proteins and subsequent kinase hyperactivity have been linked to several types of cancers (Jiang, Proc. Natl. Acad. Sci. USA 90:9026-9030, 1993; Wang, Nature 343:555-557, 1990). More recently, endogenous, highly specific protein inhibitors of cdks were found to have a major affect on cellular proliferation (Kamb et al, Science 264:436-440, 1994; Beach, Nature 336:701-704, 1993). These inhibitors include p16INK4 (an inhibitor of cdk4/D1), p21CIP1 (a general cdk inhibitor), and p27KIP1 (a specific cdk2/E inhibitor). A recent crystal structure of p27 bound to cdk2/A revealed how these proteins effectively inhibit the kinase activity through multiple interactions with the cdk complex (Pavletich, Nature 382:325-331, 1996). These proteins help to regulate the cell cycle through specific interactions with their corresponding cdk complexes. Cells deficient in these inhibitors are prone to unregulated growth and tumor formation.\nThis body of evidence has led to an intense search for small molecule inhibitors of the cdk family as an approach to cancer chemotherapy. There are no known examples of molecules related to the current invention which describe 5-substituted-indeno[1,2-c]pyrazoles as cdk inhibitors. There is one case describing indeno[1,2-c]pyrazoles having anticancer activity. There are two other examples which describe indeno[1,2-c]pyrazoles having unrelated utilities and structures.\nA series of indeno[1,2-c]pyrazoles having anticancer activity are described in JP 60130521 and JP 62099361 with the following generic structure:\nNo substitution is claimed on the indenophenyl portion of the molecule and the molecules are not indicated to be cdk inhibitors. In addition, we discovered that substitution at the 5-position was critical for cdk inhibitory activity.\nA series of indeno[1,2-c]pyrazoles having herbicidal activity are described in GB 2223946 with the following generic structure:\nThe above compounds differ from the presently claimed invention in Xn is defined as halo, alkyl, haloalkyl, and haloalkoxy; n=0-2. In addition, R1 is defined as acyl and R2 is defined as alkyl or cycloalkyl.\nA series of 1-(6′-substituted-4′-methylquinol-2′-yl)-3-methylindeno[1,2-c]pyrazoles having CNS activity are described by Quraishi, Farmaco 44:753-8, 1989 with the following generic structure:\nCompounds of this series are not considered to be part of the presently claimed invention."}
{"text": "1. Field\nExample embodiments of the following disclosure relate to an object recognition, and more particularly, to an object recognition using an analysis of point clouds of an image.\n2. Description of the Related Art\nWith respect to a method of recognizing an object, various methods are being suggested according to characteristics of sensors that are used during the process of object recognition. Among the various methods, an object recognition method that uses information of an image sensor may be largely grouped into two methods. The first method, for example, scale-invariant feature transform (SIFT) or speeded up robust features (SURF), is achieved by extracting a feature point using the brightness value of each pixel of an image, such that a descriptor is generated. The second method is achieved by an edge of an object using the differential of the brightness of an image, such that the shape of an object is inferred and recognized using the edge. The method as such is achieved by acquiring a three-dimensional (3D) point cloud using TOF (Time of Flight) and a laser sensor, in addition to an image sensor to recognize an object based on the acquired 3D point clouds.\nRegarding related methods, depending on the characteristic of the algorithm of each method, each method is provided with difficulties. As for the method, such as SIFT and SURF, of recognizing an object by the brightness value, a pattern that is not repeated on the surface of the object is needed to be provided. That is, distinguishing dishes and balls that are used in real life may be difficult. In addition, in a case when an object is largely inclined in an image, the recognition rate is rapidly decreased. The method of extracting the edge of an object, on the contrary, provides a lower object recognition rate as the external shape of the object becomes simpler. In addition, in a case when many lines are extracted due to the complexity of the background of an object, the time of performing the object recognition is rapidly increased. The algorithm using 3D point clouds as well is configured to have only the shape of an object as the characteristic of the object, and thus, the method may not be able to distinguish the objects having similar shapes to each other.\nTherefore, there is a need for improved object recognition."}
{"text": "In order to measure the surface roughness, heretofore, various methods have been used. For example, according to a stylus method, a stylus is brought into contact with a surface to be measured and moved on the surface, while according to a non-stylus or non-contact method, laser light or the like is focused to illuminate a surface to be measured and then the surface is scanned by the focused laser light to detect its reflected light. In either method, a sectional shape of the surface to be measured is obtained and a surface roughness is determined on the basis of the obtained sectional shape.\nBy the way, in order to measure the surface roughness of a work piece worked by a machine tool such as a machining center or a lathe, after the working on the machine tool has been temporarily stopped, the measurement is conducted while the work piece is removed from the machine tool and set on a surface roughness tester, or while the work piece is still mounted on the machine tool. Even with measurement of the non-stylus system, the surface roughness could not be measured in a real time manner under the circumstances using a machining liquid or cutting oil. With the measuring methods of the prior art, accordingly, it is difficult or even impossible to detect surface roughness in a real time manner during working of a work piece. Consequently, the machining has to be interrupted for the measuring process of the roughness, which has been an important reason for impeding improvement in production efficiency and product quality.\nIn order to improve the working accuracy and quality more efficiently, therefore, it has become an important task to establish the technique for clarifying surface roughness of an object to be measured (work piece) in a dynamic condition where the work piece is still being driven or worked, in other words, an in-process measurement has been required. Namely, if the measurement of surface roughness is carried out in-process manner, it becomes possible to monitor the accuracy of the surface being cut so that increase of unmanned machine tools and smooth automatic operation without troubles can be progressed, thereby greatly contributing to improvement in production efficiency and product quality.\nEven in the case of post-process measurement of the surface roughness, moreover, there is a following problem. If a narrow portion with a physical limitation such as a fine aperture impossible to insert a stylus in the stylus method, or impossible to detect reflected light in the optical method, the object has to be cut to expose the portion to be measured before the roughness is actually measured. Accordingly, it has been an important theme to establish a technique being capable of efficiently measuring the surface roughness of the narrow portion without cutting or breaking down the object to be measured.\nMain methods for measuring the surface roughness of the prior art are shown in the following patent documents 1 to 7. These are classified by measuring principles into methods based on properties of work pieces (patent documents 1 and 3), optical measuring methods (patent documents 2, 5 and 7), method utilizing styluses (patent document 4) and method using ultrasonic waves (patent document 6).    Patent document 1: Japanese Patent Application Opened No. 2003-130,606    Patent document 2: Japanese Patent Application Opened No. H11-148,812 (1999)    Patent document 3: Japanese Patent Application Opened No. 2002-340,518 (2002)    Patent document 4: Japanese Patent Application Opened No. H10-138,095 (1998)    Patent document 5: Japanese Patent Application Opened No. H6-258,060 (1994)    Patent document 6: Japanese Patent Application Opened No. H5-177,512 (1993)    Patent document 7: Japanese Patent Application Opened No. H6-99,336 (1994)"}
{"text": "FIG. 1A shows a conventional processor 100 including a plurality of cores 1050, 1051, . . . , 105N. Each of the cores 105 may include an advanced programmable interrupt controller (APIC) timer 110, a plurality of performance monitor counters 115, an instruction based sampling (IBS) fetch control register 120, and an IBS operation (Op) control register 125, which are examples of system-wide performance monitoring resources that may be simultaneously required by many software applications.\nFIG. 1B shows a conventional system including the processor 100, a plurality of drivers 1351 and 1352, a user interface 140 and a plurality of users 1451 and 1452. Referring to FIGS. 1A and 1B, a user 145 may make a time-based profile (TBP) request requiring use of an APIC timer 110, an event-based profile request requiring use of a performance monitor counter 115, or an IBS-based profile request requiring use of either an IBS fetch control register 120 or an IBS operation control register 125. However, the performance monitoring resources in each of the cores 105 of the processor 100 may not be shared by more than one user 145 at a time. Instead, once any of these performance monitoring resources are accessed by a first user, all other users are denied access to these same resources. For example, an arbitration protocol, such as one defined by Microsoft, may use a particular performance monitoring resource in the processor 100. If a first user 1451 follows the arbitration protocol, a second user 1452 will be denied access to the particular performance monitoring resource. However, if the second user does not follow the arbitration protocol, multiple accesses to the same performance monitoring resource may cause a system crash or malfunction of performance monitoring.\nThe performance monitor counters 115 in the processor 100 of FIG. 1A may be used by system software to count specific performance events associated with an operation, (e.g., retired instructions). A performance counter may be incremented in order to track the number of events that occur, (e.g., used for a sampling profile), as well as generate an interrupt when a predetermined threshold is reached, (i.e., the counter overflows). IBS is a code profiling mechanism that enables the processor 100 to select a random instruction fetch or micro-Op associated with an operation after a predetermined threshold, (e.g., number of cycles, instruction fetch count, instruction micro-Op dispatch count), has been reached, record specific performance information about the operation and generate an interrupt when the operation is completed. An interrupt handler in a driver 135 shown in FIG. 1B may then read the performance information that was logged for the operation. Instruction fetch sampling provides information about instruction translation look-aside buffer (TLB) and instruction cache behavior. The data collected for instruction fetch performance is different from the data collected for instruction execution performance.\nEvent counter multiplexing is a well-known technique used in hardware performance profiling. The event multiplexing technique allows one user to use many hardware events on a limited number of performance counters in a timeshared manner. However, by itself, event counter multiplexing does not support simultaneous performance monitoring resource usage by many other users.\nAnother well-known technique is to use a performance monitoring unit (PMU) arbitration library for resource management. However, a PMU arbitration library does not resolve conflicts to shared performance monitoring resources. For example, if a performance resource is accessed by a user, other requests to use the same performance monitoring resource are rejected.\nSince the number of available performance resources in the processor 100 is limited, the applications that need to access these resources end up competing for the limited resources. Competing needs by these applications may cause a conflict when attempting to share performance monitoring resources. Thus, a management mechanism is needed to reduce conflicts associated with sharing performance monitoring resources to avoid a system crash or malfunction of performance monitoring."}
{"text": "The present disclosure is directed to both a method and apparatus for drilling wells, and in particular for drilling wells which are drilled into geothermal formations. Geothermal wells are normally drilled to obtain steam at relatively elevated pressures and temperatures. Particular emphasis must be placed on the elevated temperature. To be sure, when an oil or gas well is drilled, there is a temperature gradient experienced as the well passes through various formations of the earth. Assume for vis h vis illustration that the surface temperature is about 75.degree. F. As a well is drilled to great depths, 10,000 feet, even 15,000 feet or deeper, there is a well known temperature gradient. Absent any geothermal activities in the near region, the temperature experienced by the well drilling equipment increases with depth at a specified rate. That rate, however, assumes the absence of a geothermal formation. A geothermal formation is characterized by the fact that heat radiates outwardly into the various geological strata between the well and the geothermal formation. The geothermal formation normally is porous having water which is captured in near proximity of the relatively hot underground formation. This fills with water which is so hot that, absent the high pressure containment, it would otherwise flash into steam. Thus, when a well is drilled into the formation, the hot water is tapped and released into the well and flashes into steam when exposed to pressure reduction sufficient to permit formation of steam. It is not uncommon for underground formations to achieve a temperature level of 500.degree. F. or 600.degree. F. or even hotter. When tapped, a typically released flow of steam, which is quite hot, travels rapidly at high velocity up the producing well to the surface. It is not uncommon to recover steam at the surface which is delivered at a temperature of perhaps about 500.degree. F. and in many instances even hotter.\nThis should be contrasted with oil production from equal depths. Ordinarily, the oil is delivered at the ambient temperature of the oil in situ in the oil bearing formation. In other words, the temperature of the oil is fairly well predictable and is normally dependent on the depths which the oil producing formation or tapped. As will be understood, the oil can sometimes be slightly hotter or colder, but it fairly well does not stray excessively from the temperature levels common for that formation.\nThe present disclosure sets forth a method and apparatus for drilling into geothermal formations which deliver extremely hot water, and more commonly which deliver saturated steam. When steam is delivered, it flows at such high velocities that the production equipment, even to the surface, is heated. Indeed, heat loss is held to a minimum because of the high velocity of the steam flow. Thus, the wellhead equipment can be installed at a high mountainous region where it is typically cold all winter and yet the wellhead equipment can quickly assume the temperature of the produced steam and maintain that temperature as long as production continues.\nOne of the uses of a geothermal well is that it provides a tremendous amount of energy in the form of heated steam for operation of steam generators and the like. The present disclosure is a method for drilling and an apparatus assisting in the production of steam from such wells. More particularly, a significant and substantial problem is encountered in the operation of such wells. Assume for purposes of illustration that the geothermal well is 6,000 ft. in depth. They can vary widely. In any event, the well typically must be cased so that the produced steam does not erode the borehole, destroy the definite shape or profile of the borehole, thereby destroying the well. The casing is installed in the well at the time of drilling. The casing string is exposed to a tremendous swing in temperature. With the swing in temperature, the casing is heated. It is heated so hot that the casing will elongate and thereby shift the wellhead equipment. Consider in a well of 6,000 ft. that the temperature differential for the casing on the average is 400.degree.. In one example, the casing is installed without steam and without heating from the producing geothermal formation at a temperature of 90.degree. along the length of the casing while after production has been initiated, the casing will be heated to 490.degree. A change of 400.degree. causes elongation of the casing sufficiently great that the casing will move the wellhead equipment unless otherwise constrained. Casing is cemented along the full length of the casing. When the change occurs, the cemented casing will be placed under a strain by virtue of the steam flow. This runs the risk of breaking the bond either between the casing and the cement or between the cement and the open hole. In either instance, it is undesirable. In both instances, it is undesirable because it may risk a migration channel along the casing on the exterior defeating the basic purpose for which the casing is normally cemented in place. As a matter of fact, cementing of the casing is normally required as a matter of good sense in well completion techniques and also as a matter of obedience to various and sundry regulatory agencies. It is highly undesirable that external migration along the cased well occur because such migration may poison clean water or artesiat formations. Artesian formations in adjacent aquafiers can be severely damaged by such migration.\nThe present disclosure enables the drilling and installation of a casing string in a geothermal well with a view of avoiding the stress which is so commonly encountered once the casing has been heated by the steam flow and a stable temperature has been established. The casing is normally heated so that the casing will move, thereby creating damage either by tearing free at one end or the other of the well, or by breaking the bond in the casing and cement structure which is formed. In the well borehole. The borehole is thus protected by the implementation of the present procedure. This procedure will be described by setting forth a number of sequential operational steps.\nIn one aspect, the present invention enables one of average skill in the art to drill and complete a geothermal well. The well is drilled by drilling through various formations of the earth toward the geothermal formation of interest. Preferably, when the well is drilled, the well is stopped short of the geothermal formation of interest. As a generalization, the depth of the geothermal formation is normally well known in advance. The depth is normally identified in a wild cat well. In a production well drilled thereafter, it is usually known where the geothermal formation will be. Alternately, with the use of temperature detectors supported in the drill string, and especially at the very bottom portion of the drill stem next to the drill bit, temperature measurements can be obtained and the temperature measurements can be used to make a prediction of when the geothermal formation will be penetrated. A geothermal formation normally is thermally isolated by the formations there above. To the extent that the geothermal formation is pervious, one or more adjacent formations tend to be impervious. This typically involves a formation which provides effective isolation to the geothermal formation so that the hot water does not penetrate the adjacent formations. While the hot water does not penetrate, heat is transferred into adjacent formations and the temperature will increase even at distances of 50 or 100 ft. from the geothermal formation targeted by the drilling process. Also, the temperature profile at depth might be well known. That adds an incremental increase in temperature as the well is drilled which can be readily known. Examples will be given below in which the, geothermal formation is located at a depth of 6000 ft. Proximity can be detected at least 200 ft. there above.\nThe procedure of the present disclosure includes the step of drilling to that specified depth. That depth will be distinguished from other depths. The well is cased at the upper portions of the well. The casing is extended to the bottom of the well drilled to that stage, and a cement shoe is then installed so that cement can be pumped down into the annular space. Separate cement jobs are implemented. The first cement job is to anchor a surface conductor string at the surface. The conductor casing string is specifically anchored by the first cement job. The surface cement job has a depth of a few hundred feet. The second cement job is pumped through the casing and cement shoe and flows into the annular space on the exterior of the casing. The second cement job will be defined below in more detail but it is primarily a cement job in which the lower end of the casing is anchored. The casing string is pulled upwardly, and the casing at the top end is cemented to anchor within the surface conductor string. The surface conductor pipe is thus cemented on both the exterior and the interior. Furthermore, the casing is left under tension, being anchored at its lower end, and is cemented in placed by the third cement job. After the casing is placed under tension, it is also equipped with a spool which surrounds the casing which thereafter grips and locks the casing and enables any surplus casing to be cut off. Wellhead equipment is completed and a Christmas tree is installed on the wellhead equipment. The latter step occurs after drilling has extended the well through the cementing shoe and into the producing geothermal formation there below.\nWhile the foregoing mentions a few of the details with regard to the procedure contemplated in practice of the present invention, the structural apparatus of the invention is likewise disclosed in some detail and is incorporated in this brief summary. The wellhead completion equipment particularly emphasizes a sleeve system which holds the casing so that it can be placed under tension."}
{"text": "1. Field of the Disclosure\nThe disclosure relates generally to the processing of data related to the subject of an insurance policy and, more particularly, to a system and method for segmenting risks at the point of application for insurance by calculating a pre-screening risk assessment score based on data received during an application process.\n2. Brief Description of Related Technology\nApplications for insurance, particularly for policies covering high-value items such as homes and other improvements to real property, typically require an assessment of risk to determine policy coverage limits of the insurer and payment rates for that coverage to be made by the prospective insured (i.e., the applicant). This risk assessment, or underwriting process, normally addresses a combination of significant sources of risk, for example, the risk that can be attributed to the property itself, and the risk that can be attributed to past behavior or other characteristics of the applicant. A physical inspection of the property often provides the measurement of risk contributed by the property itself, while an applicant's credit score often provides the measurement of risk that is likely to be contributed by the applicant.\nInspection reports have long provided valuable information in the underwriting decision-making process. Typically, a formal application process initiates the need for an inspection of the item or property to be insured. In the past, the application may have included requests for basic contact information and rudimentary details of the dwelling to determine if the property and applicant should be contacted for further detailed information or to schedule a physical property inspection. For particularly valuable items, such as a house or other real property, this inspection is typically conducted by a trained professional who travels to the property and conducts the inspection. Property inspections are typically performed by insurance agents or other designated individuals deployed in the field. Details that are physically collected about the property (e.g., the dwelling) are recorded on hardcopy or electronic forms designed to capture data related to a number of predetermined property characteristics. For instance, the inspector may utilize a number of checkboxes or placeholders on the form to identify attributes of the dwelling, such as the type of roof and siding.\nInsurers have often outsourced the process of obtaining property inspections. For example, a commercially available service provided by the assignee of the present application (under the trademarks WIDE and MAPS) routes orders for property inspections in an automated fashion to an appropriate individual (e.g., a local inspector) in the field. An inspection report is then generated in an automated fashion for electronic delivery to the insurer requesting the inspection. The service thus facilitates the front end data gathering tasks, leaving most, if not all, of the analysis of the inspection data and information for the insurer. That is, such services generally did not directly assist the insurer in the processing of the inspection data and information.\nOne exception involves a replacement cost calculation that provides an estimated cost of rebuilding any structure in the event of a total loss of property. Replacement cost calculators, or calculation services, utilize one of many industry standard models, and are available from a number of providers, including Marshall & Swift, XactValue, and e2Value. In some cases, the cost estimate is provided to the insurer with the other inspection data and information. The calculation relies on certain data or information collected via the inspection. In any event, the manner in which the input data is provided to the calculator and the manner in which the resulting cost estimate are communicated have been automated by, or within, the above-described property inspection services. For example, the calculated estimate has been displayed as part of an inspection report provided to the insurer. The insurer then compares the estimate with a coverage level of the application or current policy to determine whether an underwriting referral should be generated. In some cases, the comparison has been conducted as part of the above-described property inspection services.\nMore recently, additional capabilities have been added to the inspection data collection process to assist the insurer in other ways. For example, functionality has been developed to provide convenient access to the inspection information and data. Specifically, the MAPS™ service made available by the assignee of the present application provides data integration capabilities so that insurers can download inspection data on-line and establish computer-to-computer data interfaces. Through such interfaces, an insurer analyzes the inspection data using database tools that support, for example, searching and sorting operations. In some cases, the inspection data has also been organized via a number of categories predetermined by the insurer. More specifically, such organization results in the placement of each inspected property (or the associated inspection report) in a category, or stack, according to criteria set by the insurer. For example, each policy renewal subject to a major condition (e.g., a rotting roof) may be placed in a separate category dedicated to those policy renewals requiring a letter to be sent to the insured to identify the repairs or corrections needed for renewal. While some rules previously set by the insurer have defined one or more categories based on an action to be taken by the insurer, other rules have been set to target a number of specific issues, such as the replacement cost estimate deviating from the policy coverage by more than 5%.\nIn either case, however, such inspection reporting services provide limited to no cumulative analysis of the inspection data and information. That is, the presence or absence of various issues, a simple binary decision repeated over all of the property characteristics, merely results in a number of flags for each property. The categorization or organization of each property into stacks based on any one of those flags fails to consider the cumulative effect of them all. Such limited analysis therefore leaves a considerable amount of risk assessment analysis to a non-automated, or manual, assessment of the inspection data at the insurer, which may lead to inconsistent or incorrect underwriting referrals or other incorrect insurance coverage determinations. Moreover, as insurers consider a greater number of characteristics in these determinations, the assessment of the cumulative effect of multiple flags, i.e., issues uncovered by the inspection, will be increasingly driven by subjective judgment calls prone to error and inconsistencies.\nOne method for overcoming the shortcomings of manual inspection data assessment is described in U.S. patent application Ser. No. 11/287,649 entitled “Property Insurance Risk Assessment Processing System and Method” and hereby incorporated by reference herein. As described in the application, a method facilitates an insurance coverage determination for a property based on a plurality of characteristic elements relevant to the insurance coverage determination. The method includes (i) receiving property data indicative of attributes of the property for the plurality of characteristic elements, (ii) determining a property risk score for the property based on the property data, and (iii) supporting the insurance coverage determination with the property risk score.\nAs described above, the property inspection is an often time-consuming physical examination of the subject property. While third parties may be employed to conduct the physical inspection, and a physical inspection may provide a detailed analysis of the property to assist in the underwriting decision process, the time and cost of the physical inspection are typically necessary regardless of the final decision. In other words, the physical inspection may be a significant sunk cost that is often completed for every application before a final decision can be reached to move forward with the underwriting process.\nIn addition to property characteristics gathered by inspection, and as also previously described in U.S. patent application Ser. No. 11/287,649, the applicant's personal credit score was used as the measurement of risk contributed by the applicant. However, privacy concerns, frequent inaccuracies in the credit reporting process, and the difficulty in correcting those inaccuracies often make an applicant's credit score an illegal, unfair, or imprecise measurement of the applicant's true contribution of risk to the policy for the insurer."}
{"text": "This invention relates to a process for isomerizing alkylaromatic hydrocarbons. More particularly, this invention relates to a process for isomerizing alkylaromatic hydrocarbons with a catalytic composite comprising a combination of a platinum or a palladium component, a rhodium component, a tin component and a halogen component with a porous carrier material. The present invention utilizes a dual-function catalytic composite having both a hydrogenation-dehydrogenation function and a cracking function which affords substantial improvements in hydrocarbon isomerization processes that have traditionally used dual-function catalysts.\nProcesses for the isomerization of hydrocarbons have acquired significant importance within the petrochemical and petroleum refining industry. For example, the demand for para-xylene has created a demand for processes to isomerize other xylene isomers and ethylbenzene to produce para-xylene.\nCatalytic composites exhibiting a dual hydrogenation-dehydrogenation and cracking function are widely used in the petroleum and petrochemical industry to isomerize hydrocarbons. Such catalysts generally have a heavy metal component, e.g., metals or metallic compounds of Group V through VIII of the Periodic Table, to impart a hydrogenation-dehydrogenation function, with an acid-acting inorganic oxide to impart a cracking function. In catalysis of isomerization reactions, it is important that the catalytic composite not only catalyze the specific desired isomerization reaction by having its dual hydrogenation-dehydrogenation function correctly balanced against its cracking function, but also that the catalyst perform its desired functions well over prolonged periods of time.\nThe performance of a given catalyst in a hydrocarbon isomerization process is typically measured by the activity, selectivity, and stability of the catalyst. Activity refers to the ability of a catalyst to isomerize the hydrocarbon reactants into the corresponding isomers at a specified set of reaction conditions; selectivity refers to the percent of reactants isomerized to form the desired isomerized product and/or products; stability refers to the rate of change of the selectivity and activity of the catalyst.\nThe principal cause of instability (i.e., loss of selectivity and activity in an originally selective, active catalyst) is the formation of coke on the catalytic surface of the catalyst during the reaction. This coke is characterizable as a high molecular weight, hydrogen-deficient, carbonaceous material, typically having an atomic carbon to hydrogen ratio of about 1 or more. Thus, a problem in the hydrocarbon isomerization art is the development of more active and selective composites not sensitive to the carbonaceous materials and/or having the ability to suppress the rate of formation of these carbonaceous materials on the catalyst. A primary aim of the art is to develop a hydrocarbon isomerization process utilizing a dual-function catalyst having superior activity, selectivity and stability. In particular, it is desired to provide a process wherein hydrocarbons are isomerized without excess cracking or other decomposition reactions which lower the overall yield of the process and make it more difficult to operate."}
{"text": "U.S. Pat. No. 4,023,792 describes a sheet feeding apparatus including a storage cassette for holding a stack of paper sheets which stack is biased against some restraining or limiting member contacting the uppermost sheet. In the embodiments described, this restraining member is either pivoted out of the way or otherwise removed during the feeding process to permit the biased stack to directly contact a frictional feed roller, or the feed roller is lowered to contact the sheets. In either case the feed roller is rotated, upon contact with the paper, to urge the uppermost sheet within the stack of paper, into the paper processing path. The intermittent rotation of the feed roller while it contacts the sheets requires some form of clutch mechanism or mechanical timing arrangement to precisely control the rotation of the feed roller in order to achieve accurate timing and placement during the sheet feeding operation. In addition the constant pressure exerted on the paper stack as it is biased between the restraining member and the biasing means tends to deform the sheets and cause them to stick together so as to complicate the sheet feeding operation."}
{"text": "The inventors have studied a cross-layer optimization aimed to provide a strategy for the optimal allocation of network resources (such as sub-carriers assignments and data rates) to a set of users in consideration of system-wide fairness and spectral efficiency.\nA radio resource scheduler in a base station typically performs the allocation of subcarriers and data transmission rate to users connected to the base station for each time-slot or multiple time-slots (an epoch).\nIn a multi-user system, maximum spectral efficiency (bit/s/Hz) can be achieved by first measuring or estimating the instantaneous channel gain between each sub-carrier and each user, mapping each such value to an achievable data rate and then allocating each sub-carrier to the particular user for which the throughput (over that sub-carrier) will be maximized. When repeated for each sub-carrier in the transmission band, this simple approach maximizes the total throughput but fails to take into account the user-perceived notion of fairness, since those users who may systematically experience poor channel conditions (such as those who are close to the cell edge) will always receive fewer network resources (e.g., sub-carriers) than those who may generally experience more favorable channel conditions. On the other hand, complete fairness (whereby users are simply allocated an equal number of sub-carriers regardless of actual channel capacity) results in a loss of spectral efficiency for the entire network.\nCross-layer optimization is a tool that can be used to balance the competing interests of spectral efficiency and user-perceived fairness in the allocation of network resources. It is an optimization approach whereby the objective function of the optimization problem is a utility function.\nIn almost all wireless applications, a reliable data transmission rate, r, is the most important factor in determining the satisfaction of users. Therefore, a utility function of the data rates, U(r), should be a non-decreasing function of its argument. In particular, in a traditional network optimization, where the objective is to maximize the aggregate throughput, the utility is a linear function of the data rate, such as, U(r)=r. A utility function is typically a convex nonlinear function of the data rate or other, measurable network resources. Therefore, the utility optimization framework can be regarded as a general extension of the traditional problem of network optimization.\nThe utility functions can serve as an optimization objective for the adaptive physical and MAC layer techniques. Consequently, they can be used to optimize radio resource allocation for different applications and to build a bridge amongst the physical, medium access control (MAC) layer, and higher layers. In practice, utility functions cannot always be obtained through theoretical derivations, but may be estimated from subjective surveys. For example, for best effort traffic, a well-accepted utility function that is derived based on survey material isU(r)=0.16+0.8 ln(r−0.3)  (1)where r is in units of k b/s.\nIn order to prevent assigning too many resources to a user that already has good channel conditions, the utility function is typically selected such that its slope decreases when data rate becomes large, in accordance with the characteristic behavior of (1). Consequently, the fairness is easily taken into account.\nIn the conventional case, for which U(r)=r and the goal is to maximize the sum throughput for all users that are active at a given moment, the slope of the utility curve is a constant value of one for each user (i.e., U′(r)=1). Consequently, max-utility is achieved by assigning each sub-carrier to the user with the best channel quality, which may result in a long-term policy that unfairly allocates resources to users who typically have good channel quality and starve those users whose channel conditions are systematically less favorable (e.g., those user who are located at the cell's edge). However, the derivative of the utility function is what actually controls the threshold of comparison for the distribution of resources. Hence, a non-linear utility function whose slope decreases when the data rate becomes large may constitute a basis for a more equitable resource distribution policy.\nThe allocation and management of resources is critical for wireless networks in which a scarce collection of resources is shared by multiple users. In the current layered networking architecture, each layer is designed and operated independently. However, wireless channels suffer from time-varying multi-path fading conditions and the statistical channel characteristics of different users are different. The sub-optimality and inflexibility of this architecture result in inefficient resource use in wireless networks. Hence, there is a need for an integrated adaptive design across different layers.\nFairness and efficiency are two important issues in resource allocation for wireless networks. Traditionally, spectral efficiency is evaluated in terms of the aggregate throughput, which favors those users who are either close to the base station or who have good channel conditions. Consequently, users who are far away from the base station or who systematically have poor channel conditions may be given fewer resources. However, absolute fairness—which allocates the same quantity of system resources to all users—may lead to low spectral efficiency. Therefore, there is a tradeoff between the issues of fairness and efficiency when designing the strategies for wireless resource allocation.\nThese issues of efficient and fair resource allocation have been long studied in economics, where utility functions are used to quantify the benefit of using certain resources. Recently, utility theory has also been studied for its use in communication networks as a mechanism to evaluate the degree to which a network satisfies service requirements of users' applications. In fixed, wire-line networks, utility and pricing mechanisms have been used for flow control, congestion control and routing. In wireless networks, the pricing of uplink power control in code division multiple access (CDMA) has been investigated. Utility-based power allocation in CDMA downlinks for both voice and data applications have been proposed."}
{"text": "The present invention relates to an assembly of a buckle mechanism and blade, and a method of securing the blade to and releasing the blade from the buckle mechanism. More specifically, the present invention relates to a buckle mechanism having a locking surface, a movable guide, and a unitary biasing member.\nExisting buckle mechanisms have a moving lock member that secures a blade. These known buckle mechanisms are very complex devices comprised of a plethora of parts designed to engage and release the moving lock member from the blade to accomplish its relatively simple purpose. As a result of their complexity, known buckle mechanisms with a moving lock mechanism can be relatively expensive to manufacture and may be susceptible to failure.\nOne of the failures associated with these known buckle mechanisms is that releasing the blade from the buckle mechanism is difficult, if not impossible, after the assembly has withstood particular tensile forces. For example, vehicle passengers, after having been involved in an accident and restrained by a seat belt system employing a buckle mechanism with a moving lock, are unable to release the buckle mechanism. It is believed that the tensile forces applied to the moving lock member prevent the lock member from being moved in these conditions. Thus, a need exists for an assembly of a buckle mechanism and blade that has a minimum number of parts and avoids the difficulties of a moving lock member in a buckle mechanism.\nAccordingly, the present invention is directed to an assembly of a buckle mechanism and blade, having a fastening surface, including a base, having a surface and a projection extending away from the surface, and a guide supported by the base. The projection has a lock surface substantially perpendicular to a plane passing through the surface of the base and contiguous to the surface. The guide engages the fastening surface of the blade with the lock surface of the projection.\nIn an alternative embodiment, the buckle mechanism and blade assembly includes a base within the housing, a guide supported by the base, and at least one biasing member within the housing. The biasing member locates the guide to a first position in which the blade engages the projection and opposes movement of the guide to a second position in which the blade disengages from the projection. The base has a surface and a projection extending away from the surface, and the guide locates the blade proximate the surface.\nIn a further embodiment, the invention includes a buckle mechanism having a base with a projection and a guide supported by the base. The guide has at least one wall defining a void and is positionable in a first position so that the projection occupies at least a portion of the void and a second position so that the projection is removed from the void.\nThe present invention also provides a method of securing a blade to a buckle mechanism. In this method, the buckle mechanism has a base and a guide supported by the base, the base has a projection providing a lock surface and a portion aligned with the lock surface. The guide is configured to support the blade, and the blade has a fastening surface and a support surface adjacent the fastening surface. This method includes: (1) inserting the blade into the guide; and (2) engaging the fastening surface of the blade with the lock surface of the projection and the support surface of the blade with the portion of the base aligned with the lock surface.\nAn alternative method of the present invention is a method of releasing a blade from a buckle mechanism, where the buckle mechanism has a base and a guide supported by the base. The base has a surface and a projection extending away from the surface. The projection provides a lock surface, and the guide is configured to support the blade. The blade has a fastening surface and a support surface adjacent the fastening surface. This method includes: (1) moving the fastening surface of the blade by the guide in a first direction substantially oblique to a plane passing through the surface of the base to disengage the fastening surface of the blade from the lock surface of the projection; and (2) removing the blade from the guide in a second direction substantially parallel to the surface of the base.\nA further method of the present invention is a method of securing a blade to a buckle mechanism and releasing a blade from the buckle mechanism. The buckle mechanism has a base and a guide supported by the base. The base has a surface and a projection extending away from the surface, where the projection provides a lock surface. The guide is configured to support the blade, which has a fastening surface. This method includes: (1) inserting the blade into the guide; (2) engaging the fastening surface of the blade with the lock surface of the projection; (3) moving the fastening surface of the blade by the guide in a first direction substantially oblique to a plane passing through the surface of the base to disengage the fastening surface of the blade from the lock surface of the projection; and (4) removing the blade from the guide in a second direction substantially parallel to the surface of the base."}
{"text": "This invention relates to the sampling of blood and body fluids to inactivate or destroy infective viruses.\nThose who deal with blood and other invasively obtained body fluid samples risk infection from the samples. Those at risk include the doctor, nurse or clinical technician who takes the sample, the technicians who handle the sample and who use the sample in conducting analyses and tests, those who handle the sampling and testing equipment and apparatus, and the entire chain of individuals who attend to the disposal of sampling apparatus and the like, from the individuals who pick up the used apparatus through those who ultimately dispose of the apparatus, usually in specially designed high-temperature furnaces.\nThe risk is substantial, as evidenced by the fact that nearly all health care professionals with long experience carry the Epstein-Barr virus (EBV) and/or cytomegalovirus (CMV), the latter being probably the most ubiquitous of the pathogenic viruses. Other pathogenic viruses to which health care workers, and those who handle blood and fluid sampling and handling apparatus, are exposed include hepatitis and human immunodeficiency virus (HIV) as well as a large number of less life-threatening viruses.\nCMV is frequently associated with Pneumoncystis carinii and may cause or contribute to encephalitis and colitis and may be associated with Kaposi's sarcoma in AIDS patients. CMV carriers frequently exhibit no symptoms of infection; however, the direct and contributory effects of CMV in infectious diseases is so pervasive and subtle that a CMV infection is to be presumed if another causative agent cannot be established.\nCMV is a member of the human herpesvirus (HV) group, which are responsible for much of mankind's discomfort and pain. The herpesviruses represent a very large, clearly defined group of viruses which are responsible for, or involved in, cold sores, shingles, a venereal disease, mononucleosis, eye infections, birth defects and probably several cancers. Three subfamilies are of particular importance. The alpha subfamily includes HV-1 (herpes virus simplex 1) which causes cold sores, fever blisters, eye and brain infections, HV-2 (herpes virus simplex 2) which cause genital ulceration, and HV-3 (HV varicella zoster) which causes chicken pox, shingles and brain infections. The beta subfamily includes HV-5, the principal member of which is CMV discussed above. The gamma subfamily includes HV-4 (Epstein-Barr) which causes infectious mononucleosis and is involved in Burkitt's lymphoma and nasopharyngeal carcinoma. Additional possibly pathogenic herpes viruses no doubt exist, one type of which, HV-6, of unknown pathogenicity has been identified. (Niederman, J. C. et al., The Lancet, Oct. 8, 1988, 817.) There is evidence that the methods of this invention are effective in inhibiting the transmission of infections caused by many and perhaps all of the pathogenic herpes viruses.\nAt least one of the retroviridae is susceptible to the treatment of this invention, according to presently available data. The most notorious of the retroviridae, HIV-1, the only virus thus far identified as inducing AIDS in humans, is inactivated and/or killed using the methods and compositions of this invention. Other retroviridae are considered to be susceptible to the present invention. The enormity of the AIDS threat is so great as virtually to defy analysis. While the fear of HIV-caused AIDS is so pervasive in the United States and the more developed western world as to strike fear into all who are exposed to blood and other body fluids of a large number of individuals, the magnitude of the potential for disaster posed by HIV is, perhaps, better gauged by the fact that in central Africa and other third-world countries AIDS infects up to one-third of the population.\nHIV antibody screening of blood and organ donors has become an important but routine step in preventing the infection of the donee patient, (see Wilhelmus, K. R. et al., Ophthalmologica (Switzerland) 1987, 195(2) 57; Mal'khanov, V. B. et al., Vopr Virusol 1987, 32(2) 21; Salisbury, J. D. et al., Ophthalmic Surg 1984, 15(5) 406; Beekhuis, W. H., Doc Ophthalmol 1983, 55(1-2) 31; Collin, H. B., Arch Ophthalmol 1976, 94(10) 1726), but the risk to those who treat a large number of individuals, some of whom are certain to carry HIV, remains a very serious problem indeed.\nThe addition of detergents to various blood fractions has been described. My European Patent Specification 0 050 061, published Dec. 11, 1985, in which the term \"detergent\" is equated with the term \"amphophil\" to encompass cationic, anionic and nonionic detergents, describes the addition of various detergents to plasma protein products and suggests the addition thereof to other blood derivative products to inactivate virus and for other purposes, followed by the removal of the detergent from the product. High concentrations of detergents, from 0.25 to 10%, were required the process described in the European patent specification.\nBosslet and Hilfenhause, European Patent Specification 0 278 487, discloses that high concentrations of selected detergents inactivate certain envelope viruses.\nNeurath and Horowitz, e.g. U.S. Pat. Nos. 4,540,573, 4,481,189, and 4,591,505, indicate, however, that detergent alone is not effective as an antiviral agent in blood plasma and related products. In spite of these teachings, however, it seems safe to conclude that at least some classes of detergents in high concentrations in some types of blood derivatives do have some inactivating effect. The extent and efficacy of such procedures seems open to considerable doubt, however.\nMy copending patent application Ser. No. 07/321,522, filed Mar. 9, 1989, describes and claims a method for inactivating virus in blood samples using glycyrrhizic triterpenoid compounds.\nThe invention described and claimed herein is an improvement thereof and is based upon the discovery of unique synergistic improvement results when detergent, even trace amounts of detergent, are combined in the blood or blood product sample with glycyrrhizic triterpenoid compounds.\nThe invention described and claimed herein is an improvement thereof and is further based upon the discovery of unique synergistic improvement results when glycerol, even trace amounts of glycerol, are combined in the blood or blood product sample with glycyrrhizic triterpenoid compounds.\nFurther, the invention described and claimed herein is an improvement thereof and is further based upon the discovery of unique synergistic improvement results when glycerol, even trace amounts of ethylene diamine tetraacetic acid or salts thereof (EDTA), are combined in the blood or blood product sample with glycyrrhizic triterpenoid compounds."}
{"text": "The advent of digital media has made it possible to store large quantities of data in a very small physical footprint and has been advantageously applied in passenger entertainment systems for storing entertainment content aboard a passenger vehicle.\nDespite a resultant increase in content quality and, therefore, passenger satisfaction, however, management of digital media has become one of the biggest headaches for owners of such passenger entertainment systems. Not only has the logistics of provisioning digital media proven to be more cumbersome than traditional analog media provisioning, but the use of digital media also can result in increased operating expenses, necessitating additional capital investment. Although satellite technology and wireless network technology each can provide limited provisioning digital media capabilities, both technologies have proven to be too slow and too expensive to be practical for provisioning digital media libraries.\nThus, as usage of passenger entertainment systems continues to increase, it would be desirable to provide an affordable, effective solution for provisioning content that overcomes the obstacles of currently-available media provisioning systems.\nIt should be noted that the figures are not drawn to scale and that elements of similar structures or functions are generally represented by like reference numerals for illustrative purposes throughout the figures. It also should be noted that the figures are only intended to facilitate the description of the preferred embodiments of the present disclosure. The figures do not illustrate every aspect of the present disclosure and do not limit the scope of the disclosure."}
{"text": "Traditional “high bay” light fixtures are composed of two large spaces. One of the large spaces is formed by a reflector/refractor encasing the light source of the light fixture, and the other large space is in an electronics compartment encasing the components that control the light source. Aside from the overall large space taken up by these spaces, the electronics housing is typically situated directly above the light source. In such light fixtures, the electronic components in the electronic housing are exposed to heat generated by the electronic components themselves and to heat generated by the light sources. Even with solid state light sources, the thermal loads generated by the light source reduce the effective life of the electronics, and in turn, the effective life of the complete light fixture.\nThus, a lighting structure or fixture that reduces the damaging effects of heat generated by the light source of the lighting fixture on the electronic components of the lighting fixture may extend the effective life of the electronics and the lighting fixture."}
{"text": "Conventional on-chip transmission line structures generally have fixed impedance and fixed delay. Usually, delay and impedance cannot be arbitrarily chosen for a given transmission line. Instead, the delay and impedance are affected by the capacitance and inductance, which vary inversely to one another based upon the distance between the signal line and the ground return line(s). As such, while it is possible to change the delay of a transmission line, changing the delay comes at the cost of increasing signal loss, changing the characteristic impedance, and/or increasing the required area (e.g., footprint) of the transmission line device.\nChanging the delay of a transmission line, however, is desirable for a number of applications. For example, delay lines are utilized in signal processing operations for adjusting the time of arrival of one signal relative to that of a second signal. The delay lines may be fabricated for digital circuitry or analog circuitry, and the delay may be fixed or variable. With respect to delaying a signal having a sinusoidal waveform, this being a frequent situation in microwave applications, the effect of the delay line is to impart a phase shift; thus, in this situation the delay line may be regarded as a phase shifter.\nA plurality of phase-adjustable lines may be used in a phased array. Generally speaking, a phased array is a group of antennas in which the relative phases of the respective signals feeding the antennas are varied in such a way that the effective radiation pattern of the array is reinforced in a desired direction and suppressed in undesired directions. The relative amplitudes of, and constructive and destructive interference effects among, the signals radiated by the individual antennas determine the effective radiation pattern of the array. Phased arrays are used to electronically steer the direction of maximum sensitivity of a receiver, providing spatial selectivity or equivalently higher antenna gain. Phased arrays find use in many different wireless applications, including but not limited to RADAR and data communications. Beam steering is achieved by first shifting the phase of each received signal by progressive amounts to compensate for the successive differences amongst arrival phases. These signals are then combined, where the signals add constructively for the desired direction and destructively for other directions.\nA conventional way of controlling the phase of each element in a phased array is to provide each element with a plurality of transmission lines, each of the transmission lines having a known delay. A switch in the signal path of each element is used to select a particular transmission line for that element, thereby imparting a known delay to the element. However, such systems suffer from numerous drawbacks. For example, providing each element with a plurality of transmission lines is costly in terms of space used (e.g., footprint), manufacturing, etc. Also, the switch in the signal path of each element causes signal attenuation, which is undesirable in such applications.\nMoreover, as described above, conventional systems are incapable of changing the delay of a transmission line without either increasing signal loss, changing the characteristic impedance, and/or increasing the required area (e.g., footprint) of the transmission line device. Thus, systems that utilize delays (e.g., phased-array antenna systems) suffer from these drawbacks.\nAccordingly, there exists a need in the art to overcome the deficiencies and limitations described hereinabove."}
{"text": "1. Field of the Invention\nThe invention described herein relates to gravity-feed loading apparatus and more particularly to apparatus for delivering highly frangible nuclear fuel pellets into a container, such as a sintering boat, prior to further processing the pellets in a system.\n2. Description of the Prior Art\nThe present invention is particularly suitable for delivering green, i.e. unfired, uranium dioxide fuel pellets from the press head of a pelletizing press to a sintering boat for subsequent firing into a ceramic-like fuel pellet used in a fuel rod in a nuclear reactor.\nThe pellets have a generally cylindrical configuration, with slightly dished-in areas on opposite ends, ranging in size from about 0.3 inches to 0.8 inches long and 0.25 inches to 0.5 inches in diameter. The pellets are formed from uranium dioxide powder which is flowed into the die of the press.\nAs the power-producing capacity of the array of fuel pellets in a fuel rod is dependent upon their density, it is necessary to obtain generally uniform density among the fired, end product ceramic pellets. Such preferred density is generally on the order of 95% of the theoretical density. However, it is known that various uranium dioxide powders have sufficiently different compaction and firing characteristics so that to obtain a desired end result, the pellets must be formed in the press at different pressures depending upon the inherent characteristics of the feed powder. It is generally known that as the know-how and experience level increased in powder-producing technology, less and less pressure on feed powder in the press was required to yield the desired density of a pellet.\nThe experience of the corporate assignee of the present invention has typically conformed to the above sequence, thus, during early years of powder-forming technology, when the pressures were required to be relatively high, for example, in the order of 75,000 psi, the resultant green pellet was extremely hard and could be subjected to substantial abuse without chipping or other damage prior to firing. Therefore, the handling technique for removing fuel pellets from the press and depositing them into the firing vessel could be rather crude, and, for the most part, the pellets were permitted to tumble down an inclined ramp from the press head to above the boat and free-fall thereinto. Another arrangement, although more complex, is shown in U.S. Pat. No. 3,897,673, wherein an inclined ramp is used to deliver pellets to a loading zone. However, as pressing pressures decreased because of improved powder producing technology, the required pressing pressure became less and the green pellets formed thereby became softer, and transporting them via a continuous inclined ramp to a lower collecting point resulted in substantial chipping even when the inclined plane had a resilient surface. Such chipping was particularly evident at the end edges of the pellets. This loss of the material from the pellets had the ultimate effect of decreasing the energy producing capacity of the reactor and such damaged pellets were accordingly rejected. Scrap pellets were reworked through a process that permitted the material to be reintroduced into the powder for repressing. Due to the fact that there is always some damage to some of the pellets during manufacturing, all pellets are inspected for defects and classified accordingly as class \"A\", class \"B\", or \"scrap\". It is preferred that fuel rods be filled with class \"A\" pellets. However, due to the economics involved in reprocessing, a quantity of class \"B\" pellets (up to 10% of the quantity of class \"A\" pellets) can be used as \"bottom end\" pellets in each rod. Class \"B\" pellets that are produced in excess of the 10% limit are reprocessed along with the \"scrap\" pellets.\nThus, as the powder producing technology has advanced, it has become increasingly important to provide a mechanism for delivering the soft pellets from the pressing head to a sintering boat in a manner which minimizes damage to pellets and which requires a minimum of operator attention and adjustment of the mechanism and is relatively inexpensive and non-complex."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a device for performing a dicing in such a way that a substrate, on which two or more devices and alignment marks for positioning are formed, is positioned in accordance with the alignment mark, a positioning method, and dicing method and dicing apparatus in which the substrate is separated through the dicing to individual device elements equipped with individual devices.\n2. Description of the Related Art\nHitherto, to perform a dicing for a device substrate in which two or more devices, such as semiconductor wafers, are made, generally, there is adopted such a process that the device substrate is fixed on a dicing table through putting the device substrate on an adhesive dicing tape, the device substrate is subjected to the dicing processing so as to be segmented into device elements referred to as a chip, and the device elements are picked up one by one from the dicing tape.\nHowever, a state-of-the-art device including MEMS (Micro Electro Mechanical Systems) device, on which the research is advanced, has an extremely minute structure. Thus, when the device elements are picked up from the adhesive tape, the structure is damaged. This involves a problem that the yield is extremely low.\nIn order to improve this problem, there is proposed a method in which a device substrate is wrapped with the coagulant, without the use of the adhesive tape, and the coagulant is coagulated, so that the device substrate is fixed (cf. Japanese Patent Publication TokuKai. Hei. 10-230429, and Japanese Patent Publication TokuKai. Hei. 11-309639).\nHowever, fixing of the device substrate by the coagulation of the coagulant involves a problem of a difficulty of positioning of the device substrate. For the dicing for the device substrate, there is a need that the device substrate is exactly positioned relatively to a dicing blade. To perform the positioning, when devices are made on the device substrate, an alignment mark is formed at a position determined exactly as to the devices, and the alignment mark is photographed by a camera, so that the positioning is carried out based on the alignment mark on the image. However, when the coagulant is coagulated, the coagulation brings about a state that the coagulant is opaque or light scattering. Thus, it is difficult for a camera to recognize the alignment mark. Before the coagulation of the coagulant, it is easy for the camera to recognize the alignment mark. However, the position may be shifted in the process in which the coagulant is coagulated, and thus the recognition of the alignment mark before the coagulation would not be useful for positioning.\nIt is considered that the device substrate is fixed with a material liquid with viscous high in place of the coagulant. In this case, it is possible to recognize the alignment mark. However, as compared with a case where a position is fixed by coagulation of the coagulant, the position of the device substrate changes easily. And, the dicing for the device substrate is carried out while the coolant such as the pure water is sprayed for cooling and the removal of the cutting powder. However, even if the device substrate is covered with a liquid material of high viscosity, the device substrate is pressured by spraying of the coolant, and as a result, the device will be damaged. This involves a problem that the yield is not sufficiently improved.\nJapanese Patent Publication TokuKai. Hei. 6-331813 discloses a technology in which the substrate that forms the diffraction grating is put in the water tank that fills the pure water, and the pure water is frozen to fix the substrate and then subjected to cutting. This also easily brings about opaque or light scattering state at the time of freezing. Thus, according to the technology disclosed in Japanese Patent Publication TokuKai. Hei. 6-331813 too, it is difficult to perform a positioning base on the alignment mark.\nFurther, according to the second type of extended library apparatus show in FIG. 21, there are provided the same number of moving mechanisms as the number of coupled library apparatuses. Thus, as compared with the first type of extended library apparatus show in FIG. 19, which is capable of coping with coupling of a plurality of library apparatuses with one moving mechanism, it involves problems of high cost and increment of dissipation power."}
{"text": "1. Field of the Invention\nThe present invention relates generally to automatic production line quality control and more specifically to an automatic system for measuring various dimensions of a vehicle chassis or similar structure.\n2. Description of the Prior Art\nIn order to alleviate tedious manual work it has been previously proposed to measure the dimensions of a vehicle using a reflected light beam technique. Viz., in order to measure the various dimensions in a manner to determine the accuracy with which the vehicle chassis (or the like structure) has been assembled, the chassis is carefully disposed in a predetermined position in which a predetermined measuring coordinate grid is defined. Measuring devices which are located in predetermined locations of the coordinate grid are then used to scan the structure using laser beams or the like. The data derived from the scanning is then used to develop an image or images which can be compared with blueprint models to determine if the required degree of accuracy has been attained or not.\nHowever, this system has suffered from the drawback that even though care is taken to assure that the chassis is located in exactly the correct position within the coordinate grid, still erroneous readings tend to be produced and the need for manually executed backup checking and the like, to be executed is induced."}
{"text": "a) Field of the Invention\nThe present invention relates generally to engine retarders and, more particularly, to a system and method for supercharging or precharging an engine brake mounted to a diesel engine.\nb) Description of Related Art\nThe engine of a vehicle is often used as an auxiliary brake to retard vehicle speed. This particularly applies to heavy vehicles, such as trucks and buses, whose developments over recent years have resulted in engines of much greater power with unchanged cylinder volume. As a result, the average speed at which such vehicles are driven uphill has increased considerably, creating a need for greater braking power when driving downhill. In an exhaust brake, some form of throttle valve is incorporated in the exhaust system to achieve improved retarding power. This power, however, is relatively low and often less than half the driving power of the engine. In an engine brake, on the other hand, the exhaust valves are hydraulically opened near top dead center of the compression, thus absorbing power.\nFurthermore, the resistance of such heavy vehicles to driving has decreased over recent years, meaning that the wheel or service brakes of the vehicles are subjected to greater loads. When driving in hilly terrain, the wheel brakes should be used as little as possible, primarily for safety reasons. The average speed of the vehicle in hilly terrain is therefore greatly influenced by the available engine retarding power, which increases the requirement for a more effective retarding that will also be capable of reducing wear and tear on the wheel brakes and thereby improve running economy.\nWhen free-running, a four-stroke diesel engine, i.e. when the wheels of the vehicle drive the engine, a certain braking effect occurs as a result of the internal resistance in the engine, i.e., due to friction. This braking effect is relatively small, however, and has been further reduced in modern engines.\nA well known method of improving the engine retarding power is to mount a throttle device, for instance a butterfly valve, in the exhaust system. When the valve is closed, an overpressure is generated in the exhaust system which causes the work to increase during the exhaust stroke, with a commensurate increase in retarding power. In this manner, retarding power can be increased by placing the combustion chamber in the cylinder in communication with the exhaust system during the latter part of the exhaust stroke and during a smaller or greater part of the intake stroke.\nAlternatively, with an engine brake, the cylinder is placed into communication with the exhaust system at the end of the compression stroke and somewhat during the expansion stroke. This can be achieved, either by opening the conventional exhaust valve or with the aid of a separate valve. As a result, air compressed in the combustion chamber during the compression stroke will flow partially into the exhaust system, meaning that a large part of the compression work carried out during the compression stroke is not recovered during the expansion stroke, therewith increasing the retarding power. One known arrangement for carrying out this method utilizes the conventional exhaust valve. The exhaust valve operating cam is provided with an additional cam lobe which is operative to achieve the additional opening of the exhaust valve. The extent to which the exhaust valve is lifted by this additional lobe is relatively slight, and when the engine is used as a power source the valve clearance is sufficiently large to render the additional lobe inoperative. When braking vehicle speed with the aid of the engine, a hydraulic valve-clearance adjuster is brought into operation, such as to reduce the valve clearance, therewith bringing the additional lobe into operation The extent to which the exhaust valve is lifted during the conventional valve-opening sequence will at the same time be correspondingly greater, however, and this must be taken into account so that problems will not occur by impact of the exhaust valve against the piston.\nIt has been known that increasing boost pressure of the intake manifold increases retarding power with an engine compression release brake. Engaging a mechanical supercharger when retarding is a costly and complicated means to achieve increased boost. Likewise, adopting a supercharger with a variable area turbine is another costly and complicated means to achieve more boost.\nWhen an exhaust brake is used in conjunction with an engine brake, some of the high pressure in the exhaust manifold may be tansferred under appropriate conditions to the cylinder to provide boost pressure. However, this type of combination brake creates increased exhaust manifold pressure by restricting the airflow through the engine, resulting in increased temperature in the exhaust valve, piston and nozzle components. The increase in critical component temperature and the corresponding increase in boost pressure must be limited to avoid injector nozzle needle seat and exhaust valve seat softening and wear.\nThe need therefore exists for an improved system and method to provide engine braking without suffering from the foregoing drawbacks."}
{"text": "Raster input scanners typically employ one or more arrays such as CCDs for scanning. The array converts each scanned image line into a series of charges which, following suitable processing, are output as image signals or pixels to an end user. The scanning array, for example, may be carried on a carriage which traverses back and forth under a platen to provide the necessary relative motion between the image and array. Other arrangements such as a lamp scanning array with a movable document, etc., may also be utilized. An optical system focuses the reflected image onto the array, and one or more lamps provide illumination of the image.\nIn a typical scanning process, the image area viewed by each array of photosensors is converted into a charge potential representative of the image's gray level. Scanning takes place during an integration period of a preset duration. Following integration, the image charges are transferred to a pair of analog shift registers, the operating sequence being such that during the integration period, the image charges (image data) from a previous scanline are clocked from the shift registers, leaving the shift registers free to receive the image charges from the next integration period. The duration of the integration period, which must be sufficiently long to fully integrate the image line being scanned yet not so long as to allow the array of of photosensors to become saturated, is measured by a periodic fixed rate clock signal or shift pulses.\nWhen the scanner operates in a synchronized manner, relative scanning movement between the array and images are at a fixed rate. This in turn permits timing of the shift pulses in the signal requesting the next line of image signals, referred to as the integration signals, to be synchronized with one another.\nHowever, when a scanner operates asynchronously; i.e., when the internal transfer rate of the image input terminal is faster than the data transfer rate of the device receiving the image data; the relative scanning movement between the array and image is not fixed, but random and changes with demand. In other words, the signals are produced in response to the movement of the array or document. As a result, the timing of the integration signal timing can vary and not be in synchronism with the fixed rate shift pulses. This can reduce the integration period, resulting in an incomplete integration of the image line being scanned.\nMany methods and devices have been developed which provide asynchronous transfer when there is a difference in the transfer rate between the sending and receiving devices and also address the incomplete integration problem. These various approaches will be briefly discussed below.\nU.S. Pat. No. 4,541,061 to Schoon discloses an approach wherein the operating clock signal of a scanning apparatus is provided at a rate that matches the varying velocity of the scanning mirror through the use of a memory in which various clock rates are stored based on the known repetitive movement of the mirror. The entire contents of this U.S. Patent are hereby incorporated by reference.\nAnother approach has been proposed in U.S. Pat. No. 4,587,415 to Tsunekawa et al. This U.S. Patent discloses a photodetector with a timing control unit for controlling the information storage and information read out process of the detector. The entire contents of this U.S. Patent are hereby incorporated by reference.\nMoreover, U.S. Pat. No. 4,628,368 to Kurata et al. discloses a system for controlling the scanning rate of a document reader in which the speed, acceleration, and the deceleration of the reader are set according to image information stored in a buffer memory. The entire contents of this U.S. Patent are hereby incorporated by reference.\nU.S. Pat. No. 5,043,827 to Beikirch disclose a scanning system which corrects for incomplete integration. The scanning system utilizes an interpolator to provide a composite line of image signals when an integration period is prematurely terminated. The entire contents of this U.S. Patent are incorporated by reference\nFurthermore, U.S. Pat. No. 4,878,119 to Beikirch et al. discloses a process for operating a scanning array asynchronously. The array has at least one row of sensors for scanning an image viewed by the array during an integration period and a shift register for receiving the signal charges developed by the sensor following the integration period. This process includes periodically generating, at a constant clock rate, integration pulses defining a succession of predetermined integration intervals and generating a start integration signal in response to a demand for a line of image signals where the time at which the start integration signals occurs is different than the time at which the integration pulses occur. The process interrupts the current integration interval in response to a start integration signal to commence a new integration interval to provide the asynchronous transfer of a line of image signals. The entire contents of this U.S. Patent are hereby incorporated by reference.\nBy using asynchronous stop/stop scanning, the optical system or documents are slowed down or stopped completely until an output buffer in the image input terminal is empty enough to resume scanning. The video data timing is not synchronous to a fixed clock, but is dependent upon the motion control system. This involves a more complex control system and has motion and image quality issues associated therewith.\nIn addition to an asynchronous transfer system, a large memory buffer can be used in the image input terminal to address the problem of differing transfer rates. In this situation, a complete page can be stored in resident memory and transferred at the rate utilized by the interface of the host. However, this technique increases the cost to the system in terms of additional components call on power consumption, and the physical space required by the page memory.\nA problem exists when attempting to transfer scanned video data between an image input terminal and the host terminal when the internal data transfer rate of the image input terminal exceeds or is capable of exceeding the transfer rate of the host terminal. An example of such a situation is when the output of the scanner is gray video which is many times the volume of a threshold binary image. This produces a data rate mismatch, which can result in a truncated or corrupted page.\nAnother problem associated with transferring of data is the limited bandwidth of a synchronous clock image input terminal/host interface when used with an image input terminal which is capable of producing a wide range of data volumes and rates. An example is an image input terminal capable of producing binary and gray video data at various resolutions. In a binary mode, the data rates can be managed; however, when used in a gray output mode, the data rates become prohibited. This is especially true when the interface is governed by the host terminal which receives the data. Also, when the image input terminal produces video data which exceeds the synchronous host interface rate, such as in a gray output mode, the output video page is corrupted.\nTo prevent or avoid these problems; i.e., the mismatching of transfer rates, the present invention scans at a slower speed with the same integration time; i.e., a higher resolution, and periodically discards or skips a selected number of scanlines of valid image data to cause the scanner to realize a virtual or apparent data transfer rate which is equal to or less than the data transfer rate of the host terminal. In other words, the internal data transfer rate of the scanner is masked such that the transfer rate as perceived externally is equal to or less than the data transfer rate of the host terminal."}
{"text": "Substances which will give rise to a color change (absorption spectral change) or the like upon photoirradiation are being used for extensive development of a variety of functional materials for use in recording, display and switching.\nFor example, in high-density recording, the properties of the optical device, such as the transmittance used in the system, make it desirable to use a material which causes a change in the absorption spectrum or refractive index in the visible or near-infrared region with a short-wavelength light such as one as short as about 400 nm. In view of the increasing speed of information processing, materials are required that exhibit a very high speed response in the near-infrared region, i.e., the wavelength region for optical telecommunication. In addition, for application to electronic paper and the like, there are eagerly awaited materials which generate a color change in response not only to light but also to electric or electronic signals. However, very few materials are found which fully satisfy these requirements.\nAn ion-pair charge-transfer complex is a charge-transfer complex composed of a cationic component which serves as the electron acceptor and an anionic component which constitutes the counter ion and serves as the electron donor. The proximity of the donor and the acceptor attributable to an oxidation-reduction reaction in an ion-pair charge-transfer (sometimes abbreviated as IPCT hereinafter) complex causes a face-to-face charge-transfer interaction in addition to the coulombic interaction, and therefore the complex will develop new types of colors and properties which are not observed at all with the cation or anion alone. In the case of an ion-pair charge-transfer complex or an extended aromatic conjugated system such as X in the formula (I) set out later, photoirradiation corresponding to the absorption by the cation alone will give rise to a photoinduced electron transfer from the anion to cation thereby producing a one-electron reductant (radical cation).\nPreviously, the present inventors invented an IPCT complex of a “polymer” composed of bipyridinium cation as the backbone (Japanese Patent No. 3720277 “Polymeric photoresponsive materials and elements utilizing photoinduced electron transfer reaction”: Patent Reference No. 1). This IPCT complex polymer is suitable for use as a photo-functional material, as it causes, upon photoirradiation, an absorption spectral change (color change) at the visible to near-infrared region over a wide wavelength range for a wide range of response time. The material has the additional advantage of being easy to make into film because it is a polymer. However, the drawback is that it has a low glass-transition temperature making it difficult to apply the material to a device.    Patent Reference No. 1: Japanese Patent No. 3720277 “Polymeric photoresponsive materials and photoresponsive elements utilizing photoinduced electron transfer reaction”."}
{"text": "An electronic component such as a heat-sensitive resistive element is used, as a resistor, for information communication devices such as moving body communication terminals or personal computers, as well as electronic devices such as wearable devices, medical devices, consumer devices, or automobile electrical devices.\nRecently, a reduction in a thickness of such electronic devices has been required, and development of electronic components having a small thickness has been under development due to limitation of a thickness dimension of the electronic devices.\nFor example, Patent Literature 1 describes a thin film thermistor using a ceramic substrate having a thickness dimension of 50 μm to 300 μm in a chip resistor to be mounted on a board. Patent Literature 2 describes a chip resistor having a thickness dimension of 60 μm to 150 μm. On the other hand, Patent Literature 3 mentions a problem that, when a ceramic substrate in which an insulating substrate has a thickness dimension of 30 μm to 100 μm is used, the substrate may break at the time of fabricating a chip component. Therefore, even when a chip resistor having a thickness dimension of 60 μm according to Patent Literature 2 can be formed, there is concern that the substrate may break, a problem such as cracks at the time of mounting of the board occurring, and a product with extremely low reliability being obtained.\nDue to this, it is preferable to use an extremely thin substrate having a thickness dimension of 100 μm or less so that the substrate does not break at the time of fabricating a chip component, and to fabricate a chip resistor and a heat-sensitive resistor using an insulating substrate formed of, for example, a ceramic material of which the reliability such as the stability or heat resistance is high.\nOn the other hand, in a resistor and a heat-sensitive resistor that are used for a wearable device that monitors biological information, a catheter that is a medical device, or the like, it is required for the resistor and the heat-sensitive resistor to be formed of a material for which biocompatibility is taken into consideration regarding a material to be used when there is a high risk of exposure of the resistor and the heat-sensitive resistor inside a living body.\nWhen a part of the heat-sensitive resistor is formed of a biocompatible material, direct contact with a living body is possible, and an effect that accurate temperature detection of the living body is possible can be expected."}
{"text": "Air fresheners are known for use in motor vehicle passenger cabins. One often seen variety includes the air freshener “trees” that may be observed hanging from a vehicle inside mirror.\nThe air freshening agent may be provided as a solid, liquid or gel material. In the case of liquids, the liquid may be provided in a bottle with a wick to volatize the freshening agent into the air, alternately liquids may be saturated into or absorbed into sponge, cloth or foam carrier materials from which components of the freshener liquid volatize or evaporate into the air.\nU.S. Pat. No. 5,407,642 discloses a clip mounted air freshener assembly having a holder into which an air freshening element is inserted.\nU.S. Pat. No. 6,123,906 discloses a unitary fork-shaped air freshener of impregnated polypropylene material and insertable into an air vent discharge.\nU.S. Pat. No. 6,1976,263 discloses an electrically heated air freshener for automobiles, wherein the fragrance dispensing unit includes electrical heating elements heating replaceable gel-scent cartridges and an electrical plug for receiving electrical power from a cigarette lighter socket.\nIt is also known to be advantageous to provide air filter elements (also known as cabin air filters) arranged within and integrated within the passenger compartment air ventilation duct system of a motor vehicle to remove various types of particulate contaminants. While some vehicles are manufactured and equipped with cabin air filters, a greater number of motor vehicles lack provisions for or to install a cabin air filter in the motor vehicle. In such vehicles it is at least difficult and maybe nearly impossible to retrofit a cabin air filter into the air ventilation duct system to provide cabin air filtration.\nTherefore there remains a need in the art for an air freshening clip cabin air cleaning device including an air filter device with an integrated air freshening agent which is low in cost and easily installed into or onto the air ventilation outlet grill, which is adapted to filter at least a portion of the air ventilation duct airflow while providing an aromatic air freshening agent to diffuse a pleasant scent into the air."}
{"text": "(1) Field of the Invention\nThe present invention relates to a joint positioning device, and more particularly, to an adjustable joint positioning device.\n(2) Description of the Prior Art\nThere are many joints in our body to connect bones and the joints allow the limbs to be bent to take different exercises. However, the joints are easily injured because they are used so frequently and the tendons are easily injured. Once the joints hurt, the operation of the bones are affected and the patients are restricted to move because the injured joints have to be well positioned.\nA conventional joint positioning device is shown in FIG. 8 and generally includes a top support frame 70 and a bottom support frame 80. The top and bottom support frames 70, 80 are composed of two restriction frames 71, 81 and each restriction frame includes two straight sections 72/82 on two sides thereof and a transverse section 73/83 which is connected between the two straight sections 72/82 and are curved so as to match the curvature of the patient's thigh or calf. Soft pads 75, 85 and straps 76, 86 are connected to the inside of the top and bottom support frames 70, 80. The straps 76, 86 are tied to the thigh and calf to achieve the purpose of restriction of the joint.\nThe restriction frames 71, 81 each are an integral part so that the curvatures of the transverse sections 73, 83 have to match the size of the thigh and calf of the patients. Therefore, the manufacturers prepare the restriction frames 71, 81 with different curvatures of the transverse sections 73, 83. A high stocking cost is a pressure to the manufacturers so that only three or four sizes of the transverse sections 73, 83 are produced which obviously cannot meet the needs of patients. If the curvature of the transverse sections 73, 83 is not matched with the patients' thigh and/or calf, the patients feel uncomfortable and the result for wearing the joint positioning devices is not satisfied.\nIn addition, the restriction frames 71, 81 are integral parts and cannot be folded which occupy a large space and increase the cost of transportation.\nThe present invention intends to provide a joint positioning device wherein the curvature can be adjustable to match the patients of different sizes."}
{"text": "With a conventional radiographic apparatus, an X-ray source projects X-rays to an object (e.g., a medical patient). X-ray beams transmitted through the object are detected to radiograph the object by a screen film cassette, film autochanger, CR (Computed Radiography), FPD (Flat Panel Detector), or the like.\nIn the field of radiography, a high-resolution solid X-ray detector which uses an FPD is proposed. This solid X-ray detector has an X-ray sensor including a two-dimensional array in which arrays of 3,000 to 4,000 photoelectric conversion devices (e.g., photodiodes) are arrayed two-dimensionally. Each photoelectric conversion device generates an electrical signal corresponding to the X-ray dose projected to the X-ray sensor. The X-ray image of an object is obtained by arranging the object between an X-ray source and X-ray sensor and converting the X-ray dose which has been transmitted through the object into an electrical signal. A signal from each photoelectric conversion device is read out individually and digitized, and thereafter image-processed, stored, or displayed.\nFIG. 9 is a conceptual view showing the structure of a system which includes a conventional cassette type radiographic apparatus. As shown in FIG. 9, a radiographic apparatus 801 incorporates an X-ray detector 802. X-rays generated by an X-ray generator 803 irradiate an object 804. The X-rays transmitted through the object 804 are detected by photoelectric conversion devices (not shown) which are arrayed like a matrix on the X-ray detector 802. Electrical signals output from the photoelectric conversion devices are image-processed by an image processor 805 to display the X-ray image of the object 804 on a display 806 such as a monitor.\nIn recent years, a low-profile, higher-density mounting technique has improved, and a compact, low-profile solid X-ray detector which uses an FPD is becoming possible, e.g., an X-ray screen film cassette (see Japanese Patent Laid-Open Nos. 2003-057352 and 2002-186604).\nFIG. 10 is a side sectional view showing an example of a cassette type radiographic apparatus which uses an FPD. As shown in FIG. 10, an electronic cassette used for radiography or the like has a particle phosphor 131, e.g., GOS, which converts X-rays into visible light, a MIS photosensor portion 115 which uses amorphous silicon and photoelectric conversion devices 109 which are arranged like a matrix to convert the visible light into electrical signals, a TFT switching portion 116, a base 110, a circuit board 111 which supports the base 110, a circuit board 113 on which electronic components for processing the photoelectrically converted electrical signals are mounted, wiring lines 114, case lid 101 which is used to house the above members, and a case main body 117. A buffer member 102 serving as a relaxation portion that relaxes a force from outside the housing is arranged between the case lid 101 and particle phosphor 131. A resin 130 made of PET or the like is arranged between the buffer member 102 and particle phosphor 131.\nThe particle phosphor 131 is adhered to the photoelectric conversion devices 109 by an adhesion layer 106 through a second protection layer 107 made of an organic substance such as PI and a first protection layer 108 made of a nitride or the like. The circuit board 111 on which the electronic components 113 for processing the electrical signals photoelectrically converted by the photoelectric conversion devices 109 are mounted is planarly attached to the lower surface of the base 110 in tight contact through an insulating sheet 112.\nConventionally, in the electronic cassette which uses the particle phosphor 131 such as GOS, the photoelectric conversion devices 109 which are made of a component such as glass are more vulnerable to an external force than the particle phosphor 131 and protection layers 107 and 108. Therefore, a cassette type radiographic apparatus which uses an FPD has been designed with reference to the strength of the photoelectric conversion devices 109.\nIn the conventional cassette type radiographic apparatus which uses the FPD, problems occur when a columnar crystal phosphor such as CsI (cesium iodide crystal) is used in place of the particle phosphor 131. The reasons of the problems are roughly classified into two. According to the first reason, the columnar crystal phosphor is fractured by a weaker external force than the photoelectric conversion devices made of a component such as glass. As the columnar crystal phosphor is more vulnerable to the external force, the conventional technique for protecting the photoelectrical conversion devices cannot sufficiently protect the columnar crystal phosphor. According to the second reason, as the columnar crystal phosphor has stricter demands for an external force acting on a small area than a granular phosphor used in the conventional cassette type radiographic apparatus using the FPD, the crystals of the columnar phosphor may be broken. When a stress is applied, holes may be formed in the protection film by steps formed by variations in crystal length.\nAccording to the prior art, a standing- or lying-position radiographic apparatus is available which uses a columnar crystal phosphor and an FPD. When this apparatus is used as a portable cassette type radiographic apparatus, problems occur. The reasons of the problems are roughly classified into two. According to the first reason, unlike the standing- or lying-position portable cassette type radiographic apparatus, the portable cassette type radiographic apparatus is used in various applications, and accordingly the weight of the object may be applied to it. Therefore, a stress absorbing portion which is not necessary in the standing- or lying-position radiographic apparatus is necessary in the portable cassette type radiographic apparatus in case various types of external forces are applied when, e.g., the operator erroneously hits the apparatus with something from above or places his or her elbow on the apparatus. According to the second reason, the cassette type radiographic apparatus is not provided with a grid space. The standing- or lying-position radiographic apparatus usually uses a scattered ray removing mechanism (grid), and the scattered ray removing mechanism (grid) serves as a stress absorbing portion. In the cassette type radiographic apparatus, however, the radiographic target is often a body portion, e.g., a limb, which does produce many scattered rays, and accordingly a scattered ray removing mechanism (grid) is not usually mounted due to the requirements for a lower profile and lighter weight. Therefore, a structure is necessary which protects the columnar crystal phosphor from the external force.\nIn the cassette type radiographic apparatus, for achieving a lower profile., a gap must be minimized as small as possible in the direction of thickness of the cassette. With a small gap, however, if something is erroneously dropped on the case to deform it elastically, the phosphor may be broken at a high possibility. In particular, in the case of a columnar crystal phosphor such as CsI, when the stress acts on the crystals to break the phosphor crystals, the directional characteristics of light in the phosphor change, and a scar may be undesirably present in an image obtained by X-ray irradiation.\nFIGS. 11A to 11D show images obtained when a stress acts on the columnar crystal phosphor to break the phosphor crystals. After an experiment (FIG. 11A) of dropping a screwdriver onto the columnar crystal phosphor (CsI: TI+) was conducted, the columnar crystal phosphor was adhered to photoelectric conversion devices, and radiography was performed (FIG. 11B). As shown in FIG. 11B, a large hitting mark ranging for a diameter of 3 mm appears on the image. FIG. 11C is a graph showing the sectional profile of an output value of a photodetector taken along the line A–A′ of FIG. 11B. As shown in FIG. 11C, a portion where the optical output increases by about 10% and a portion where the optical output decreases by about 10% are present. FIG. 11D Is a graph showing the sectional profile of the output value of the photodetector taken along the line B–B′ of FIG. 11B. As shown in FIG. 11D, many portions are obviously present where the optical output increases by about 10%. In this manner, with reference to FIGS. 11C and 11D, portions where the output value Increases or decreases are obviously present, and a difference in output value is as large as ±10%. The output value fluctuates probably because the phosphor crystals are broken and a light-outputting portion changes. As is apparent from the experimental images of. FIGS. 11A to 11D, since the columnar crystal phosphor has a low strength, even when only a lightweight material such as a screwdriver is erroneously dropped, it forms a scar In the image to cause a clinical problem."}
{"text": "1. Field of the Invention\nThis invention relates to a nickel-base alloy which resists corrosion, abrasion, erosion, cavitation and temperature cycles and possesses a high-temperature strength. In particular, this invention relates to an improved alloy of the type disclosed in corresponding application Ser. No. 828,718 which alloy has the following composition in percent by weight:\nCarbon: 0.5 to 5 PA1 Tungsten: 2 to 15 PA1 Chromium: 25 to 55 PA1 Boron: 0.5 to 3.5 PA1 Silicon: 1 to 5 PA1 Iron: 1 to 5 PA1 Nickel: balance PA1 (a) percent by weight C=percent by weight C.sub.1 +percent by weight C.sub.2 +0.2; PA1 (b) percent by weight W=36.times.percent by weight C.sub.1 ; and PA1 (c) percent by weight Cr.gtoreq.17+9.times.percent by weight-4.times. PA1 percent by weight C.sub.1 =carbon content required to form carbides with refractory metals (tungsten, molybdenum) PA1 percent by weight C.sub.2 =carbon content required to form chromium carbide PA1 0.2=a constant which is equivalent to the solubility limit of C in the nickel-chromium solid solution.\nand meet the following boundary conditions:\nWhere:\nIn particular this invention relates to such an alloy additonally containing phosphorus, e.g., in an amount of 0.5-4.5 weight percent. The alloy has improved processing characteristics.\n2. Discussion of the Prior Art\nIn virtually all fields of industrial activity, it is desired to solve problems relating to the wear and corrosion of highly stressed components by coating them with hard alloys so as to save raw material. Hard alloys based on cobalt and nickel have been used for decades with success particularly in the manufacture of machines, equipment and fittings.\nIn nuclear engineering the use of cobalt alloys, so-called stellites, is limited by the presence of cobalt. For this reason it has been attempted most recently to replace stellites by nickel-chromium-boron-silicon hard alloys, which can also be processed in a simpler manner. On the other hand, it has been found that boride- and silicide-containing hard nickel alloys are not satisfactory with respect to their resistance to corrosion, particularly intergranular corrosion. In many cases, nickel-chromium-boron-silicon alloys are unsatisfactory also in respect of high-temperature strength and susceptibility to cracking. For this reason, a new type of nickel-containing hard alloy has been developed and disclosed in corresponding application Ser. No 828,718, filed Aug. 29, 1977, entitled: \"Nickel-Base Alloy\", assigned to the assignee hereof, the disclosure of which is hereby incorporated specifically herein by reference.\nIn view of their characteristic content of hard material, these coating alloys are described as carbide-content of the alloy in adaptation to the relatively high chromium and tungsten contents and additions of boron, silicon and iron permit of a production of hard alloys having a hardness in the range of 35 to 50 to HRC and a high resistance to corrosion.\nThe new hard aloys have the structure which is typical of hard alloys. The hard material phases which are responsible for the wear resistance are embedded in a regular distribution in a relatively ductile metal matrix. In dependence on the alloy constituents and their mutual solubility in the solid state, the metal matrix consists of a solid solution of nickel, which contains chromium, tungsten, silicon, carbon and iron. The predominant hard material phase consists of a complex carbide of the type M.sub.7 C.sub.3, wherein M stands for the metals chromium, nickel and tungsten. Owing to the extremely low solubility of boron in the matrix solid solution, another hard phase consists of chromium borides, mostly of the type CrB. Depending on the proportions of the alloying constituents within the stated tolerance ranges and the cooling conditions adopted in the production and processing of the alloy, additional hard phases may be stabilized in small quantities, if desired. They consist of the complex carbides of types M.sub.23 C.sub.6 and M.sub.2 C and a silicon-stabilized carbide of the type M.sub.6 C.\nThe new alloys afford advantages regarding the resistance to abrasion at higher temperatures because alloys having a comparatively smaller hardness at room temperature can be used so that the processing and machining will be facilitated.\nIn addition to its high strength at room temperature, the alloy distinguishes by an excellent strength at high temperature and an almost constant strength at temperatures in the range of 300.degree. to 600.degree. C.\nAs regards thermal expansion and density, the novel hard nickel alloys are comparable to corresponding conventional hard alloys based on nickel or cobalt.\nIn addition to their good mechanical properties, carbide-containing hard nickel alloys exhibit also a favorable melting behavior. The melting temperature range of about 1170.degree. to 1230.degree. C. is lower than that of comparable hard cobalt alloys. This results in excellent processing characteristics.\nOwing to their excellent resistance to corrosion, the novel hard alloys close a gap which has existed in the conventional hard nickel alloys, which have an extremely poor resistance particularly to oxidizing media."}
{"text": "Examples of conventionally known self-ballasted fluorescent lamps include a self-ballasted fluorescent lamp which is provided with a cover, a lighting circuit contained in the cover, and an arc tube bent or otherwise formed into an appropriate shape and contained in a globe, said cover having a base that can be mounted in a socket designed for a typical light bulb.\nA self-ballasted fluorescent lamp commercially available at present typically has such specifications as a height of approximately 130 mm (including the height of the base), an outer diameter of approximately 70 mm, an outer tube diameter of the arc tube of approximately 12 mm, a discharge path length of approximately 280 mm, a tube wall thickness of not less than 1.1 mm, and a lamp power of approximately 13 W. Due to its configuration, however, it is difficult to provide a fluorescent lamp which has such an arc tube and is as compact as typical light bulb. Nevertheless, there is an increasing demand for fluorescent lamps made to an even smaller scale.\nAnother example of self-ballasted fluorescent lamps is disclosed in Japanese Patent Laid-open No. 1987-12051, which relates to a fluorescent lamp, wherein an arc tube having three U-shaped bent bulbs is disposed in such a way that the three U-shaped bent bulbs respectively correspond to the three sides of an approximately equilateral triangle. However, as there is no detailed discussion in the above Japanese Patent Public Disclosure as to various criteria regarding the reduction of the dimensions of the lamp, such as dimensions and the shape of the arc tube as well as criteria for lighting the lamp, the invention disclosed in said publication does not provide the optimum configuration for reducing the dimensions of the lamp.\nAnother example of fluorescent lamps is disclosed in Japanese Patent Laid-open No. 1987-12051, wherein the arc tube of the fluorescent lamp is formed in a U-like shape having corners bent at approximately 90°. However, the configuration having such an arc tube, i.e. an arc tube having sharp corners, presents a problem of irregularity in luminance, because the corners of the arc tube are too close to the globe when the arc tube is contained in the globe which is as small as that of a typical light bulb.\nAnother example of fluorescent lamps is disclosed in Japanese Patent Laid-open No. 1997-69309, wherein the arc tube is bent into a spiral or other shape so as to produce a lamp having a shape and dimensions nearly identical to those of a typical light bulb. However, a configuration which calls for bending the arc tube into such a complicated shape as a spiral requires a complicated production process and presents a problem in that reduction of production costs is difficult. As it is difficult to put such an arc tube in practical use for reasons described above, an arc tube having U-shaped bent bulbs is normally used. However, a lamp having such an arc tube, too, is difficult to be made compact, because it imposes various limitations in the shape and the dimensions of the U-shaped bent bulbs.\nWhen the dimensions of a fluorescent lamp are reduced, there arises the danger of heat from the arc tube exerting an unfavorable influence on the lighting circuit that is contained in the cover. As a fluorescent lamp disclosed in Japanese Patent Laid-open No. 1996-273615, one of the known ways to solve this problem is a configuration which calls for disposing a circuit board for mounting components of the lighting circuit thereon in such a manner that the components are positioned apart from the ends of the arc tube at which the electrodes are provided. As a result of the reduction of the dimensions of fluorescent lamps, however, circuit boards, too, are made compact. Therefore, the above configuration presents a problem in that the reduction in the space in which the necessary components are mounted increases the planar dimensions of the lamp too much, particularly at the part where the cover is located.\nRegarding a self-ballasted fluorescent lamp which is provided with a cover having a base that can be mounted in a socket designed for an incandescent lamp, a lighting circuit contained in the cover, and an arc tube bent or otherwise formed into an appropriate shape and contained in a globe, a configuration which calls for disposing a circuit board at the base-facing end of an arc tube that is bent in a U-like shape and arranging electrical components on both end of the circuit board is widely known. One of examples of such configuration is disclosed in Japanese Patent Laid-open No. 1988-245803. Compared with the aforementioned configuration which calls for positioning the circuit board apart from the ends of the arc tube, said configuration disclosed in Japanese Patent Laid-open No. 1988-245803 is more effective in reducing the horizontal dimensions of the lamp at the region of the cover. On the other hand, it presents such problems that interference between the electrical components and the arc tube, especially between the electrical components and the end of the arc tube, increase the influence of heat exerted on the electrical components and that such a configuration makes the lamp too long.\nAs described above, the outer diameter of the conventional self-ballasted fluorescent lamp has larger than the outer diameter of typical light bulb. Therefore, this configuration presents problem in that it is not suitable for a luminaire which uses a typical light bulb in place of a typical light bulb.\nRegarding a fluorescent lamp which is used to a self-ballased fluorescent lamp, a configuration which a bulb has a bent discharge path, which formed by connecting a three U-shaped tubular bodies in series and electrodes disposed at the both ends of the bulb is widely known. One of examples of such configuration is disclosed in Japanese Patent Laid-open No. 220360-1989. And such configuration ensures the length of a discharge path and the reduction of the dimensions of a fluorescent lamp.\nIn some cases, such a fluorescent lamp uses a main amalgam for controlling the pressure of the mercury vapor in the bulb within an appropriate range during the time that the lamp is lit under normal conditions and an auxiliary amalgam for absorbing mercury floating in the bulb when the lamp is turned off and releasing the absorbed mercury during the early stage of lighting, including the moment when the lighting is initiated. In a configuration where the amalgams are used, the main amalgam is contained in a minute tube which serves to discharge the air and projects from an end, i.e. the end at which an electrode is contained in the bulb, of a tubular body that is located at an end of the bulb, while the auxiliary amalgam is disposed at an end of a tubular body positioned at the middle portion of the bulb. However, when the exhaust minute tube that contains the main amalgam and projects from a tubular body situated at an end of the bulb is located at the same end at which an electrode enclosed in the bulb is located, the temperature of the main amalgam becomes too high due to the influence of the heat from the electrode. Such an increase in the temperature of the main amalgam impairs the effective control of the pressure of the mercury vapor and causes the pressure of the mercury vapor to increase too much, resulting in a decrease in luminous flux. This configuration presents another problem in that it is difficult to uniform or stabilize the pressure of the mercury vapor in the tubular body that is located at the other end of the bulb, at a long distance from the main amalgam.\nAnother example of fluorescent lamps characterized by inclusion of a main amalgam is disclosed in Japanese Utility Model Publication No. 1992-47893, wherein a main amalgam is disposed in a minute tube projecting from an end of one of the three tubular bodies that form the bulb, said tubular body being the middle tubular body of the three tubular bodies. The fluorescent lamp having this configuration is capable of reducing the influence of the heat from the electrodes exerted on the main amalgam, limiting the pressure of the mercury vapor within an appropriate range by preventing an excessive increase in temperature of the main amalgam, and also capable of reducing the distances from the main amalgam to the respective ends of the bulb by a nearly identical degree, thereby making the pressure of the mercury vapor uniform and stable throughout the interior of the bulb. On the other hand, the above configuration presents a problem in that disposing the main amalgam in the minute tube projecting from an end of the middle tubular body of the three tubular bodies of the bulb positions the main amalgam too far from the electrodes, making it difficult to warm the main amalgam. Especially at the initiation of lighting, when both the ambient temperature around the fluorescent lamp and the temperature of the main amalgam itself are low, the main amalgam is slow to release mercury, because it takes a long time for the temperature of the main amalgam to reach the level where the main amalgam functions most effectively. As a result, the luminous flux build-up characteristics become poor, and it takes an excessively long time to stabilize the luminous flux.\nIn response to the recent tendency toward compact fluorescent lamps, the demands for reduction of the dimensions of bulbs are on the increase. In the configuration where each minute tube for discharging the air is provided at an end of the bulb, the reduction of the diameter of the bulb makes it necessary to reduce the diameter of the minute tubes. However, a minute tube having a diameter smaller than a given dimension has poor exhaust conductance, resulting in decrease in the exhaust efficiency. On the other hand, if the diameter of the minute tubes are not reduced, the distance between each minute tube and a pair of inner copper-weld wires that support an electrode is reduced, making the operation of sealing the bulb difficult.\nIn order to solve the above problems, an object of the present invention is to provide a fluorescent lamp and a self-ballasted fluorescent lamp that are characterized by the minutes tubes positioned so as to permit the amalgams to be disposed at the optimum locations and enable the reduction of the diameter of the bulb. Another object of the present invention is to provide a self-ballasted fluorescent lamp and a luminaire that enable the reduction of the diameter of the portion near the base of the lamp as well as the reduction of the length of the lamp. Yet another object of the present invention is to provide a compact self-ballasted fluorescent lamp which is equivalent to a typical light bulb and a luminaire including such a self-ballasted fluorescent lamp."}
{"text": "1. Field of the Invention\nThe present invention relates to an image recording apparatus for recording images on a recording member, and in particular relates to a conveying device to be provided in the apparatus for conveying a member to be conveyed such as the recording member.\n2. Description of the Related Art\nAn inkjet recording apparatus having an electrostatic-attraction-type conveying device for conveying a member to be conveyed such as a recording member will be described as a conventional image recording apparatus of such a kind, for example, with reference to FIGS. 6 to 7B. FIG. 6 is a schematic sectional view of a conventional image recording apparatus, and FIGS. 7A and 7B are a plan view and an enlarged sectional view of a conveying device for a recording member, respectively.\nA recording apparatus 50 shown in FIG. 6 comprises a recording head 51 for recording images and a conveying belt 52, which is an electrostatic-attraction-type conveying member.\nThe recording head 51 is of a line type in that a number of nozzles are arranged in rows to have substantially the same width as that of the recording member, and can perform rapid image-forming in comparison with a serial type in which a recording head mounted on a carriage is scanned.\nHowever, because of the high speed, ejected ink is unfixed just after recording, so that if anything is brought into contact with a recorded surface, images may be marred. Such a phenomenon may more or less occur even in the serial type image recording apparatus.\nThen, in such an apparatus, the rear surface of the recording member is held with the conveying belt 52 having electrostatically-attracting means so as to convey the recording member without coming in contact with anything.\nThe conveying belt 52 is made by connecting both ends of a band member to be annular, and stretched and looped around a driving roller 58 and a plurality of supporting rollers 59 to 61, so that the recording member such as a paper sheet or plastic sheet is conveyed by rotating the conveying belt 52.\nAs the electrostatically-attracting means of the conveying belt 52, there is an electrode group (referred to as interdigital electrodes 53 and 54 below), in which strip-shaped electrodes with different polarities are alternately arranged as shown in FIG. 7A, and there are conductive brushes 55 arranged on both sides of the conveying belt 52 as power feeding means, as shown in FIG. 7B.\nIn such a structure, when one of the interdigital electrodes 53 and 54 is grounded while a predetermined voltage is applied to the other so as to produce a difference in potential, an attracting force due to electrostatic power can be obtained.\nHowever, the conventional technique described above has the following problems.\nThe conveying belt 52 is driven by the driving roller 58 at a high speed, and foreign particles such as paper dust and oil are stuck to the driving roller 58 during the long time running, so that the frictional force between the driving roller 58 and the conveying belt 52 may be reduced or jamming due to plugged paper or dust may be produced, resulting in the abnormal driving of the conveying belt 52.\nIf such a case is neglected, the abnormal belt driving may finally result in a serious problem such as head damage or electric leakage, so that first hand stopping has been required.\nThen, in order to detect the abnormal belt driving of the conveying belt 52, an encoder 56 attached to the supporting roller 59 is detected with a photosensor 57. Alternatively, although not shown, the driving state of the conveying belt is detected by reading an encoder pattern printed on the conveying belt with a photosensor.\nHowever, as described above, the conventional conveying belt has had additional procedures and cost for attaching an encoder to the conveying belt or printing an encoder pattern thereon in order to detect the driving state of the conveying belt."}
{"text": "Carbon blacks may be utilized as pigments, fillers, reinforcing agents and for a variety of other applications. Carbon blacks are widely utilized as fillers and reinforcing pigments in the compounding and preparation of rubber compositions and plastic compositions. Carbon blacks are generally characterized on the basis of their properties including, but not limited to, their, surface areas, surface chemistry, aggregate sizes and particle sizes. The properties of carbon blacks are analytically determined by tests known to the art, including iodine adsorption surface area (I.sub.2 No.), nitrogen adsorption surface area (N.sub.2 SA), dibutyl phthalate adsorption (DBP), dibutyl phthalate adsorption of the crushed carbon black (CDBP), cetyl-trimethyl ammonium bromide absorption value (CTAB), Tint value (TINT), Dmode and .DELTA.D50. It is advantageous in the compounding and preparation of rubber and plastic compositions to utilize a carbon black that is easily dispersed in the media.\nCarbon blacks are also widely utilized as pigments in the formulation of ink compositions, paints and the like, wherein it is generally desirable to use a carbon black pigment which can be easily dispersed. For example, newsink compositions are made in two stages. First the carbon black pigment and a dispersing vehicle, comprising resin, additives, and oil or solvent, are mixed to form a \"premix\" and then the premix is ground, for example, in a shot mill, to complete the dispersion of the carbon black in the ink composition. Dispersion of the carbon black pigment in the ink composition occurs during the formation of the premix and during the grinding of the premix. A carbon black pigment which is easy to disperse will allow an ink maker to produce an ink in a reduced period of time, which results in improved economy. For the above reasons, and others, it would be advantageous to produce an easily dispersible carbon black pigment.\nCarbon blacks are generally produced in a furnace-type reactor by pyrolyzing a hydrocarbon feedstock with hot combustion gases to produce combustion products containing particulate carbon black. A variety of methods for producing carbon blacks are generally known.\nIn one type of a furnace carbon black reactor, such as shown in U.S. Pat. No. 3,401,020 to Kester et al., or U.S. Pat. No. 2,785,964 to Pollock, hereinafter \"Kester\" and \"Pollock\" respectively, a fuel, preferably hydrocarbonaceous, and an oxidant, preferably air, are injected into a first zone and react to form hot combustion gases. A hydrocarbon feedstock in either gaseous, vapor or liquid form is also injected into the first zone whereupon pyrolysis of the hydrocarbon feedstock commences. In this instance, pyrolysis refers to the thermal decomposition of a hydrocarbon. The resulting combustion gas mixture, in which pyrolysis is occurring, then passes into a reaction zone where completion of the carbon black forming reaction occurs.\nIn another type of a furnace carbon black reactor a liquid or gaseous fuel is reacted with an oxidant, preferably air, in the first zone to form hot combustion gases. These hot combustion gases pass from the first zone, downstream through the reactor, into a reaction zone and beyond. To produce carbon blacks, a hydrocarbonaceous feedstock is injected at one or more points into the path of the hot combustion gas stream. The hydrocarbonaceous feedstock may be liquid, gas or vapor, and may be the same or different than the fuel utilized to form the combustion gas stream. Generally the hydrocarbonaceous feedstock is a hydrocarbon oil or natural gas, however other hydrocarbonaceous feedstocks such as acetylene are known in the art. The first (or combustion) zone and the reaction zone may be divided by a choke or zone of restricted diameter which is smaller in cross section than the combustion zone or the reaction zone. The feedstock may be injected into the path of the hot combustion gases upstream of, downstream of, and/or in the restricted diameter zone. Furnace carbon black reactors of this type are generally described in U.S. Pat. Reissue No. 28,974 and U.S. Pat. No. 3,922,335 the disclosure of each being incorporated herein by reference.\nIn generally known reactors and processes, the hot combustion gases are at a temperature sufficient to effect pyrolysis of the hydrocarbonaceous feedstock injected into the combustion gas stream. In one type of reactor, such as disclosed in Kester, feedstock is injected, at one or more points, into the same zone where combustion gases are being formed. In other type reactors or processes the injection of the feedstock occurs, at one or more points, after the combustion gas stream has been formed. The mixture of feedstock and combustion gases in which pyrolysis is occurring is hereinafter referred to, throughout the application, as \"the effluent\". The residence time of the effluent in the reaction zone of the reactor is sufficient, and under conditions suitable, to allow the formation of carbon blacks. In either type of reactor, since the hot combustion gas stream is continuously flowing downstream through the reactor, pyrolysis continuously occurs as the mixture of feedstock and combustion gases passes through the reaction zone. After carbon blacks having the desired properties are formed, the temperature of the effluent is lowered to a temperature such that pyrolysis is stopped, thereby halting the further production of carbon blacks.\nIn generally known processes the lowering of the temperature of the effluent to stop pyrolysis is accomplished by injecting a quenching fluid, through a quench, into the effluent. As generally known to those of ordinary skill in the art, pyrolysis is stopped when the desired carbon black products have been produced in the reactor. One way of determining when pyrolysis should be stopped is by sampling the effluent and measuring its toluene extract level. Toluene extract level is measured by ASTM D1618-83 \"Carbon Black Extractables--Toluene Discoloration\". The quench is generally located at the point where the toluene extract level of the effluent reaches an acceptable level for the desired carbon black product being produced in the reactor. After pyrolysis is stopped, the resulting mixture of combustion gases and carbon black generally passes through a heat exchanger to further cool the mixture. This heat exchanger is often advantageously utilized to preheat the combustion air to be utilized in the process while at the same time cooling the quenched mixture from the reactor. Thus, this heat exchanger is often referred to as the combustion air heat exchanger.\nAfter passing through the combustion air heat exchanger the quenched mixture passes through a secondary cooler to further cool the mixture. A secondary quench is typically utilized for the secondary cooler. The purpose of the secondary cooler is to further lower the temperature of the quenched mixture to a temperature such that the bag filter system utilized to separate the carbon blacks will not be damaged.\nAfter further cooling of the mixture by the secondary cooler the cooled mixture passes downstream into separating means whereby the carbon blacks are recovered. The separation of the carbon black from the gas stream is readily accomplished by conventional means such as a precipitator, cyclone separator or bag filter. This separation may be followed by pelletizing using, for example, a wet pelletizer.\nThe temperature of the effluent in the reactor is generally above about 1750.degree. F., often reaching over 3300.degree. F. In conventional processes, the quench that stops pyrolysis cools the effluent to below about 1650.degree. F., often to about 1400.degree. F. The resulting quenched mixture enters the heat exchanger (combustion air heat exchanger) at this temperature and is further cooled by the heat exchanger to about 1000.degree. F. The secondary cooler further cools the mixture to about 500.degree. F., a temperature that will generally not damage the separation means such as the bag filter system.\nAs set forth above, in heretofore generally utilized carbon black production processes, pyrolysis of the carbon black yielding feedstock in the hot combustion gas stream is stopped by a quench injecting a quenching fluid, typically water. The use of a water quench may disadvantageously result in tiny drops of water being carried with the quenched mixture into the combustion air heat exchanger contacting the heat transfer surfaces thus causing a hard fouling layer of carbon black to build up on the internal heat transfer surfaces of the heat exchanger. This hard fouling layer of carbon black is difficult to remove, and decreases heat transfer through the surface, and thus is disadvantageous.\nAlso, occasionally in heretofore generally utilized carbon black production processes bits of the hard fouling layer are broken loose and entrained into the product stream due to large changes in the operating conditions of the reactor. These bits of hard fouling layer pass through the reactor and end up mixed in with the collected carbon black product. The material mixed in the carbon black product makes the product difficult to disperse.\nAdditionally the water droplet evaporation process, of the quench water, causes the formation of micropellets of carbon blacks in the gas stream. These micropellets are more difficult to disperse in the final application media."}
{"text": "This application is based on and claims priority under 35 U.S.C. xc2xa7 119 with respect to Japanese Application No. 2000-080969 filed on Mar. 22, 2000, the entire content of which is incorporated herein by reference.\nThe present invention is directed to a vehicular brake control device which makes it possible to apply a brake fluid pressure to each of wheel cylinders irrespective of a brake pedal operation.\nOne of the prior art vehicular brake control devices of the kind, as is well known, includes a hydraulic pressure generation means which makes it possible to cause a master cylinder to generate a master cylinder pressure in response to a brake pedal depression force which is, for example, boosted by a booster, a hydraulic pressure supply means which supplies the master cylinder, as a brake fluid pressure, to a wheel brake cylinder of each of wheels, and a pressure application means which introduces a hydraulic pressure into a pressure application chamber in the hydraulic pressure supply means. In addition, the aforementioned known vehicular brake control device includes a control means which controls the pressure application means to automatically adjust the master cylinder pressure and which drives the hydraulic pressure supply means to control a braking force to be applied to each of the wheels. The hydraulic pressure generation means has a normally opened electromagnetic valve, as a pressure-maintaining valve, which establishes and interrupts a fluid communication between the master cylinder and each of the wheel cylinder and a normally closed electromagnet valve, as a pressure-decrease valve, which establishes and interrupts a fluid communication between the master cylinder and a reservoir.\nIn the foregoing conventional or prior art brake control device, when the master cylinder pressure is automatically increased, if the wheel cylinders are for front-wheels as driven wheels of an FR-vehicle under traction control, the pressure retaining valve is required to close which is provided between the master cylinder and each of such wheel cylinders in order to prevent an increase of the brake fluid pressure of each of the wheel cylinders.\nHowever, if the pressure retaining valve malfunctions or has trouble remaining in its opened state, the wheel cylinder which is never increased in hydraulic pressure is pressurized. Thus, for example, despite traction control for starting or accelerating the vehicle, a problem occurs such as prevention of starting the vehicle due to the brake fluid application into the wheel cylinder, or so-called xe2x80x98brake dragxe2x80x99.\nThus, a need exists to provide a vehicular brake control device, for overcoming the aforementioned problem, which make it possible to detect a malfunction of an electromagnetic valve wherein its opened state remains unchanged or can not be shifted to its closed state.\nThe present invention has been developed to satisfy the request noted above and a first aspect of the present invention is to provide a vehicular brake control device which comprises:\na master cylinder which is capable of generating a brake fluid pressure in response to a brake pedal depression force;\nhydraulic pressure control means for controlling a brake fluid pressure which is transmitted from the master cylinder to each of wheel cylinders of respective wheels;\npressure application means which is capable of automatically applying a hydraulic pressure to each of the wheel cylinders independent of brake pedal depression by way of the hydraulic pressure control means; and\ncontrol means for controlling a braking force of each of the wheels by driving the pressure application means and/or the hydraulic pressure control means,\nthe hydraulic pressure control means including a first electromagnetic valve which establishes and interrupts fluid communication between the master cylinder and each of the wheel cylinders and a second electromagnetic valve which establishes and interrupts fluid communication between a reservoir and each of the wheel cylinders,\nthe control means including pressure application check means which makes all the first electromagnetic valves of the respective wheels and all the second electromagnetic valves of the respective wheels closed and opened, respectively, at a predetermined timing and which initiates a pressure application check for judging whether or not automatic pressure application is made successfully by controlling the pressure application means,\nthe pressure application check means having an electromagnetic valve abnormal condition judging means which judges that the first electromagnetic valve of one of the wheel cylinders fails to be kept opened if a changing rate of the hydraulic pressure applied by the pressure application means increases after a fixed duration (KT1) has elapsed from the initiation of the pressure application check.\nIn accordance with a first aspect of the present invention, the first electromagnetic valve and the second electromagnetic valve are closed and opened, respectively, if a failure or malfunction is found in closing the first electromagnetic valve (i.e., the first electromagnetic valve fails to shift from its open position to its closes position), the brake fluid flows into the reservoir through the second electromagnetic valve, the consumed amount of the brake fluid at each of the wheel cylinders increases. Thus, the increasing speed of the hydraulic pressure becomes slower which is introduced by the pressure application means after initiation of the pressure application check, which results in that immediately when the reservoir is filled with the brake fluid the increasing speed of the hydraulic pressure changes drastically, which is introduced by the pressure application means. Therefore, if a speed change is found to be increased in the increasing speed of the hydraulic pressure which is introduced by the pressure application means within a fixed time duration (KT1), measured from the initiation of the pressure application, the electromagnetic valve abnormal judging means makes it possible to detect, with great accuracy, a failure or malfunction in closing the first electromagnetic valve.\nA second aspect of the present invention is to provide a vehicular brake control device, as a limited version of the first aspect, wherein the pressure application check means judges, if the master cylinder pressure increases to a predetermined value (KP1) before a fixed time duration (KT2) elapses after the initiation of the pressure application check, that the pressure application means and the first electromagnetic valve of each of the wheels are in normal condition.\nIn accordance with the second aspect of the present invention, if the hydraulic pressure which is introduced by the pressure application is increased to a fixed value (KP1) until a fixed time duration (KT2) elapses after initiation of the pressure application check, it becomes possible to judge that the pressure application means and the first electromagnetic valve of each of the wheels are in good order or in normal condition.\nA third aspect of the present invention is to provide a vehicular brake control device, as a limited version of the first aspect, wherein the pressure application check means includes pressure application abnormal judging means and pump regenerative voltage check means, the pressure application abnormal judging means judging temporarily that the pressure application means is abnormal if more than a fixed time duration (KT2) from the initiation of the pressure application check is required for the master cylinder pressure to increase to a predetermined valve (KP1), the pump regenerative voltage check means driving a pump for a short time duration in order to suck the brake fluid, after the temporal judgment, which enters the reservoir by way of the second electromagnetic valve of each of the wheel cylinders, the pump regenerative voltage check means judging that if the regenerative voltage of the pump after its driving is less than a predetermined value (KVP) the first electromagnetic valve of one of the wheel cylinders fails to be kept open.\nIn accordance with the third aspect of the present invention, if the master cylinder pressure requires a fixed time duration (KT2) measured from the initiation of the pressure application check for reaching a predetermined value (KP1), it is impossible to judge whether the pressure application means is abnormal or the first electromagnetic valve is in failure in its closing operation. Thus, the pressure application means is temporarily judged to be abnormal. Thereafter, if the pump regenerative voltage is found to be not greater than the predetermined value (KVP) after driving the pump for a short time duration, the brake fluid flows into the reservoir by way of the second electromagnetic valve due to a failure or malfunction in the first electromagnetic valve, which makes it possible to judge that the pump inertia force of the pump upon its termination is used to discharge the brake fluid and therefore the pump regenerative voltage lowers temporarily. Thus, in such a case, it can be judged that a failure or malfunction is found in closing the first electromagnetic valve of any one of the wheel cylinders.\nA fourth aspect of the present invention is to provide a vehicular brake control device, as a limited version of the third aspect, wherein the pump regenerative voltage check means judges that if the regenerative voltage of the pump after its driving is not less than the predetermined value (KVP), each of the first electromagnetic valves is normal and the pressure application means is abnormal.\nIn accordance with the fourth aspect of the present invention, if the regenerative voltage of the pump after its driving is not less than a fixed value (KVP) the first electromagnetic valve is normal and the brake fluid fails to flow into the reservoir by way of the second electromagnetic valve. Thus, the foregoing pump inertia force is not used to discharge the brake fluid when the pump stops, i.e., the pump has no load, which makes it possible to judge that the first electromagnetic valve is normal and the pressure application means is abnormal."}
{"text": "Existing techniques for offloading data traffic and data service flows from a Wireless Wide Area Network (WWAN) (e.g., a 3GPP-LTE or 802.16m (WiMAX) network) to a Wireless Local Area Network (WLAN) (e.g., a Wireless Fidelity (WiFi) 802.11a/b/g/n network) include the static mapping of services to either WWAN or WLAN links. Some network-based offloading techniques, such as those specified in various 3GPP standards, require the client to choose whether to receive a particular service over either a WLAN or WWAN link without coordination between the WLAN and WWAN networks. Additionally, these techniques use IP layer offloading and require interworking between 3GPP and WLAN networks, impacting 3GPP core network elements such as Authentication, Authorization, and Accounting (AAA) servers, the Home Subscriber Server (HSS), policy components, and the like.\nThus, there are general needs for improved methods of coordinating user and control plane solutions for transporting service flows across the WLAN and WWAN networks. There are also general needs for improved methods of controlling WLAN offload based on relevant network conditions and metrics of interest."}
{"text": "Ever since the first patent for this type of instrument was issued to an American, Mr. Hayen (U.S. Pat. No. 3,947,968) in 1874, special publications of this nature have appeared time and again, and a search of international Patent Gazettes issued during the last decade or so reveals that available ellipsograph designs fall into five categories:\n(1) The string type utilises the principle that the sum of the distances of any point on the ellipse to two focal points is constant.\n(2) The projection image copying type is designed on the principle that all kinds of ellipses can generally be drawn on paper along the edge of the shadow mimicked by a circular board model inclined at various degrees in space. This category may be typified by U.S. Pat. No. 4,194,292.\n(3) The obliquely cut cylinder type is prepared by ensleeving another leg on one of the legs of the dividers, enabling it to rotate and slide up and down with the drawing pen, forming an adjustable included angle to it. This is typified by U.S. Pat. No. 4,204,327. Ellipsographs designed according to the above three principles cannot describe ellipses that are very small, nor those that are fairly elongated.\n(4) The planetary gear type is typified by the ellipsograph disclosed in U.S. Pat. No. 4,226,022. It can hardly overcome the crawling which causes lines drawn in this way to have a saw-tooth appearance.\n(5) The cross groove type includes the majority of ellipsographs designed like those disclosed in DT2743352, SU454139, and BG1754588. This type has the advantage of being able to draw elongated ellipses while its defects include the following: it is difficult to ensure the matching precision of the slide groove and the coincidence of the elliptical line drawn in opposite directions; its grooved frame tends to be clumsy; it requires that the unevenness of the paper surface"}
{"text": "1. Field of the Invention\nThe present invention relates to medical devices for monitoring vital signs, e.g., arterial blood pressure.\n2. Description of the Related Art\nFalse alarms generated by conventional vital sign monitors can represent up to 90% of all alarms in critical and peri-operative care, and are therefore a source of concern. A variety of factors cause false alarms, one of which is motion-related artifacts. Ultimately false alarms can have a severe impact on the safety of hospitalized patients: they can desensitize medical professionals toward ‘true positive’ alarms, lead them to set dangerously wide alarm thresholds, or even drive them to completely disable alarms. This can have a particularly profound impact in lower-acuity areas of the hospital, i.e. areas outside the intensive care unit (ICU), emergency department (ED), or operating room (OR), where the ratio of medical professionals to patients can be relatively low. In these areas a single medical professional (e.g. a nurse) often has to care for a large number of patients, and necessarily relies on automated alarms operating on vital sign monitors to effectively monitor their patients.\nStudies in critical care environments indicate that the majority of false positive alarms are simple ‘threshold alarms’, meaning they are generated when a patient's vital sign exceeds a predetermined threshold. Patient motion can result in a vital sign having an erroneous high or low value, which in turn can trigger the false alarm. In most cases, these alarms lack any real clinical meaning, and go away after about 20 seconds when they are not acknowledged. Alarms can also be artificially induced when a patient is moved or manipulated, or if there is an actual problem with the vital sign monitor. False alarms due to motion-related artifacts are particularly very high when measured from ambulatory patients.\nBlood pressure is a vital sign that is particularly susceptible to false alarms. In critical care environments like the ICU and OR, blood pressure can be continuously monitored with an arterial catheter inserted in the patient's radial or femoral artery. Alternatively, blood pressure can be measured intermittently using a pressured cuff and a technique called oscillometry. A vital sign monitor performs both the catheter and cuff-based measurements of blood pressure.\nAlternatively, blood pressure can be monitored continuously with a technique called pulse transit time (PTT), defined as the transit time for a pressure pulse launched by a heartbeat in a patient's arterial system. PTT has been shown in a number of studies to correlate to systolic (SYS), diastolic (DIA), and mean (MAP) blood pressures. In these studies, PTT is typically measured with a conventional vital signs monitor that includes separate modules to determine both an electrocardiogram (ECG) and pulse oximetry (SpO2). During a PTT measurement, multiple electrodes typically attach to a patient's chest to determine a time-dependent ECG component characterized by a sharp spike called the ‘QRS complex’. The QRS complex indicates an initial depolarization of ventricles within the heart and, informally, marks the beginning of the heartbeat and a pressure pulse that follows. SpO2 is typically measured with a bandage or clothespin-shaped sensor that attaches to a patient's finger, and includes optical systems operating in both the red and infrared spectral regions. A photodetector measures radiation emitted from the optical systems that transmits through the patient's finger. Other body sites, e.g., the ear, forehead, and nose, can also be used in place of the finger. During a measurement, a microprocessor analyses both red and infrared radiation detected by the photodetector to determine the patient's blood oxygen saturation level and a time-dependent waveform called a photoplethysmograph (‘PPG’). Time-dependent features of the PPG indicate both pulse rate and a volumetric absorbance change in an underlying artery caused by the propagating pressure pulse.\nTypical PTT measurements determine the time separating a maximum point on the QRS complex (indicating the peak of ventricular depolarization) and a foot of the optical waveform (indicating the beginning the pressure pulse). PTT depends primarily on arterial compliance, the propagation distance of the pressure pulse (which is closely approximated by the patient's arm length), and blood pressure. To account for patient-dependent properties, such as arterial compliance, PTT-based measurements of blood pressure are typically ‘calibrated’ using a conventional blood pressure cuff and oscillometry. Typically during the calibration process the blood pressure cuff is applied to the patient, used to make one or more blood pressure measurements, and then left on the patient. Going forward, the calibration blood pressure measurements are used, along with a change in PTT, to continuously measure the patient's blood pressure (defined herein as 'cNIBP). PTT typically relates inversely to blood pressure, i.e., a decrease in PTT indicates an increase in blood pressure.\nA number of issued U.S. Patents describe the relationship between PTT and blood pressure. For example, U.S. Pat. Nos. 5,316,008; 5,857,975; 5,865,755; and 5,649,543 each describe an apparatus that includes conventional sensors that measure an ECG and PPG, which are then processed to determine PTT. U.S. Pat. No. 5,964,701 describes a finger-ring sensor that includes an optical system for detecting a PPG, and an accelerometer for detecting motion."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an apparatus for stacking and soldering individual package devices for a three dimensional stack package device, and more particularly to such a stacking and soldering apparatus allowing automatic stacking and soldering by applying solder flux to metal portions of the individual package devices during the transfer of the package devices.\n2. Description of the Related Arts\nIn present packaging technology, a three dimensional stacking of several integrated circuit devices has attracted attention in order to meet demands of high density and increased capacity. The three dimensional stack package device can be obtained by stacking non-packaged semiconductor bare chips or by stacking packaged individual devices. When the packaged devices are stacked, most of the existing assembly processes such as a die attachment, a wire bonding and a molding can be applied to the individual packages prior to stacking. However, specific additional processes are necessary in stacking the individual package devices.\nIn the package stacking technology, it is an important factor how to electrically interconnect the stacked devices. A soldering method is one of the widely used methods for the electrical interconnection.\nFIG. 1A is a schematic diagram of a conventional soldering apparatus which is disclosed in U.S. Pat. No. 5,236,117, and FIG. 1B is a detailed view of FIG. 1A. Package devices 12 to be stacked are applied with a flux, and then fixed to the left end of an impact arm 10. When the metal part, e.g., outer leads of the package devices are heated to a predetermined temperature, the packages 12 are dipped into a molten solder fountain 8. The impact arm 10 is pivotally mounted on pivot assembly 6, which allows the left end of the impact arm holding the package device to move upward and downward. The impact arm 10a in an elevated position and the impact arm 10b in the lowered position are shown in FIG. 1A. When the impact arm 10 moves to its lower position, the outer leads of the package are dipped into the molten solder in the fountain 8. The height stopper 2 is to set the lowered position of the impact arm. Vibration means 4 is to vibrate the impact arm 10 in order to prevent solder voids or lack of solder tinning of closely spaced leads.\nWith this dipping technique, the solder application can be performed more efficiently. However, when the lead pitch gets finer as the package device requires more and more I/O pins, a solder bridge and therefore electrical shorting actions can occur. For solving these problems, impact weight 9 drops onto the top of the opposite end of the impact arm 10. Therefore, the package device and the impact arm are rapidly returned to the upper position, so that all of the excess molten solder is removed from the leads. The movement of the impact weight 9 is controlled by an air cylinder 7 which in turn may be controlled by a programmable controller.\nIn the conventional soldering apparatus explained above, loading process of the vertically stacked individual packages to the apparatus is performed manually. On the other hand, J-lead type individual packages such as SOJ (Small Outline J-bend), TSOJ (Thin SOJ) and PLCC (Plastic Leaded Chip Carrier) are suitably applied to the soldering for vertical interconnection of the three dimensional stack package device. Thus, if an automated stacking and soldering apparatus is provided for the exclusive use for the J-lead type packages, the mass production of the stack package device can be possible."}
{"text": "(1) Field of the Invention\nThe present invention relates to a method for manufacturing hydrochlorofluoroolefins. More particularly, this invention relates to synthesis of 2-chloro-3,3,3,-trifluoropropene (HCFC-1233xf).\n(2) Description of Related Art\nFluorocarbon based fluids have found widespread use in industry in a number of applications, including as refrigerants, aerosol propellants, blowing agents, heat transfer media, and gaseous dielectrics. Because of the suspected environmental problems associated with the use of some of these fluids, including the relatively high global warming potentials associated therewith, it is desirable to use fluids having the lowest possible greenhouse warming potential in addition to zero ozone depletion potential. Thus there is considerable interest in developing environmentally friendlier materials for the applications mentioned above.\nTetrafluoropropenes, having essentially zero ozone depletion and low global warming potential, have been identified as potentially filling this need. However, the toxicity, boiling point, and other physical properties in this class of chemicals vary greatly, even between different isomers of a compound. One tetrafluoropropene having valuable properties is 2,3,3,3-tetrafluoropropene (HFO-1234yf). HFO-1234yf has been found to be an effective refrigerant, heat transfer medium, propellant, foaming agent, blowing agent, gaseous dielectric, sterilant carrier, polymerization medium, particulate removal fluid, carrier fluid, buffing abrasive agent, displacement drying agent and power cycle working fluid. Thus, there is a need for new manufacturing processes for the production of tetrafluoropropenes and in particular 2,3,3,3-tetrafluoropropene.\nOne process for producing tetrafluoropropenes involves the use of tetrachloropropenes as a reactant (US 2007-0197842 A1). Additionally, several other methods of preparing hydrofluoroalkenes are known. For example, U.S. Pat. No. 4,900,874 (Ihara et al) describes a method of making fluorine containing olefins by contacting hydrogen gas with fluorinated alcohols. U.S. Pat. No. 2,931,840 (Marquis) describes a method of making fluorine containing olefins by pyrolysis of methyl chloride and tetrafluoroethylene or chlorodifluoromethane. The preparation of HFO-1234yf from trifluoroacetylacetone and sulfur tetrafluoride has been described. See Banks, et al., Journal of Fluorine Chemistry, Vol. 82, Iss. 2, p. 171-174 (1997). Also, U.S. Pat. No. 5,162,594 (Krespan) discloses a process wherein tetrafluoroethylene is reacted with another fluorinated ethylene in the liquid phase to produce a polyfluoroolefin product.\nNotwithstanding prior teachings, applicants have come to appreciate a continuing need for methods of efficiently preparing intermediates of certain hydrohalocarbons, particularly compounds which are in part useful as intermediates in the preparation of tetrafluoropropenes, such as 2,3,3,3-tetrafluoropropene (HFO-1234yf).\nThe prior art discloses various processes for the preparation of polyhaloolefin products that involve separate steps as well as disparate reaction conditions, reagents, and catalysts. The efficiency of such multi-step processes is thus limited by the efficiency of each individual step. As such, one inefficient step may make the entire process more resource intensive, less effective at converting intermediates to the desired fluorocarbon products and less productive, suffering yield losses due to increased impurity formation.\nThe present invention offers a less resource-intensive process that produces increased conversion of intermediates to the end product polyhaloolefin over a longer period due to substantially increased catalyst life. This is achieved by Applicants' discovery of subjecting a tetrachloropropene to fluorination with a fluorinating agent and a catalyst in the presence of a stabilizer."}
{"text": "Although the present invention is useful in any adhesive composite needing a delivery system, it has particular benefit in connection with transparent film dressings and surgical drapes. These dressings and drapes are widely used as a protective layer over a wound, facilitating healing in a moist environment while acting as a barrier to liquids and bacteria. Dressings of this type are available under trade names such as Tegaderm.TM. (3M, St. Paul, Minn.), Bioclusive.TM. (Johnson & Johnson, New Brunswick, N.J.), Op-Site.TM. (T. J. Smith & Nephew, Hull, England) and Uniflex.TM. (Howmedica, Largo, Fla.).\nThe polymeric films used in such dressings are conformable. By this it is meant that these films are extremely thin, flimsy, and supple. They are supplied with a releasable protective liner overlying the adhesive coated surface of the film. When the liner is removed, the adhesive coated film tends to wrinkle and stick to itself interfering with the smooth aseptic application of the dressing or drape to the skin. Various delivery systems have been proposed to obviate this problem. PG,4\nOne such delivery system utilizes a removable frame to support the edges of the film during its application. U.S. Pat. No. 4,372,303 describes a composite having a relatively thin polymeric film which is conformable to animal anatomic surfaces. A pressure-sensitive adhesive is coated on one surface of the polymeric film, a release liner is attached to the adhesive coated surface of the film, and a frame is attached to the exposed surface of the film. In use the frame stays with the film until the film is affixed to the skin, at which time it is removed.\nAnother delivery system is described in U.S. Pat. No. 4,513,739. That patent describes a wound dressing comprising a film that is coated on one face with a biocompatible adhesive. One or more liner sheets are releasably adhered to the adhesive coated surface. Release retarding means are provided along one edge or a pair of opposed edges of the dressing to require a greater force to separate the layers (film, adhesive, and liner) at the edge or edges than at the remainder of the contact area. The release retarding means may be a thicker strip of the polymeric film, a stiffer piece of film adhered to the film at its edges, a strip of liner wrapped around the edge of the film and adhered to the exposed surface of the film, or a strip attached to the exposed surfaces of both the film and the liner. Alternatively, a more aggressive adhesive may be used at the edges than on the central portion of the dressing.\nA third delivery system is described in U.S. Pat. No. 4,485,809. In this system, a window frame of liner material is provided for support during placement of the dressing by removal of a center section of the release liner. The window frame is then removed by separation of the film at the perforations provided within the area of the frame.\nA fourth delivery system is used in the Veni-gard.RTM. dressing (Conmed, Utica, N.Y.). In this dressing, a permanently attached foam frame is provided on the non-adhesive side of the film backing.\nAn additional problem which occurs in the use of thin film adhesive dressings is the premature release of the dressing from the skin, generally initiated by failure of the adhesive at the edges of the dressing. Practitioners will occasionally secure the edges of the dressing with medical tape in order to avoid this problem. This is a very inefficient solution to the \"edge lift\" problem involving additional materials and time."}
{"text": "The preferred consumer digital video interfaces are High Definition Multimedia Interface (HDMI) and Digital Visual Interface (DVI). DVI is commonly used by PC equipment to drive digital video displays. DVI typically supports 24-bit RGB at video rates up to 165 MHz. A DVI driver accepts 24-bit RGB data and serializes it into three serial channels. The video clock is added as a fourth channel. As the RGB data is serialized, it is encoded using an 8b/10b encoding scheme called Transition Minimized Differential Signaling (TMDS). HDMI is backwards compatible with DVI. It supports alternate (non-RGB) color spaces and includes the ability to carry digital audio.\nAs with DVI, HDMI data is encoded to represent active video periods and control periods. In addition, HDMI includes a third entity called a Data Island. Data Islands are used to communicate additional data during the blanking interval. For example, data islands are used to send digital audio data in HDMI.\nData delivered using a DVI or HDMI interface may be encrypted using High-Bandwidth Digital-Content Protection (HDCP). Implementation of HDCP requires a set of unique secret device keys. During authentication, the receiver will only receive content once it demonstrates knowledge of the keys. Furthermore, to prevent eavesdropping and stealing of the data, the transmitter and receiver will generate a shared secret value that is consistently checked throughout the transmission. Once authentication is established, the transmitter encrypts the data and sends it to the receiver for decryption.\nHeretofore, DVI/HDMI data has only been deliverable using hard wires due at least in part to the way the data source and the display device need to communicate with each other. However, it may be desirable for a consumer to want to place a display device (such as a flat-panel television) on a wall opposite from the video source (e.g., receiver or DVD player). In this case, the consumer would typically be required to purchase and install DVI cabling from the video source to the display. However, this may be both costly and present difficult installation issues.\nA wireless radio frequency (RF) system could be used between a video source and a display device. However, there may be numerous reasons for preferring a wireless optical configuration between a digital video source and a display device. For example, the hardware required for an optical wireless signal may be less complex than a RF system. Moreover, a wireless optical solution is more secure since it will not penetrate walls as with an RF system.\nTherefore, a wireless optical system which eliminates the need for cabling between a digital video source and a digital video sink may be desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to a child carrier apparatus and its operating method.\n2. Description of the Related Art\nAn infant stroller can be typically used for transporting an infant in a convenient manner. Because the adult pushes the infant stroller from the rear side, the child seated in the infant stroller generally cannot see the adult. This may be disadvantageous as the sight of a familiar face may comfort the child. To overcome this issue, certain manufacturers have developed infant strollers provided with a backrest that can be adjustable toward the rear or front of the stroller for either forward or rearward configurations of use. More specifically, when the backrest is inclined forward, the child can be seated facing the rear of the stroller. When the backrest is inclined rearward, the child can be seated facing the front of the stroller. Because this approach requires a seat of a larger size, the adult may erroneously place the child on the front or rear portion of the seat while the stroller is not configured to the adequate orientation. For example, the adult may have the child seat on the rear portion of the seat while the stroller is in the forward configuration of use, or reciprocally. Such misuse may cause overload, and affect the safety of the stroller.\nTherefore, there is a need for an improved design that is convenient to operate and can address at least the foregoing safety issues."}
{"text": "The present invention relates generally to a system and method for tracking game of chance play and proceeds. More particularly, the invention relates to a system and method for allocating a portion of a player\"\"s purchase and/or potential winnings to a player\"\"s account.\nGames of chance (e.g., lotteries) have been offered by various sponsors (e.g., government agencies and private industry) depending upon the laws of each State. With the proliferation of States allowing casinos, river boat gambling and lotteries, increasingly, sponsors have engaged in competition for the gaming public\"\"s money. For example, State-operated games of chance (e.g., lotteries) are engaged in competition with other forms of legal gaming (e.g., church-sponsored bingo, casinos, online gaming and/or racing) and illegal gaming. At present, there are approximately thirty-seven States that operate lotteries and games of chance to benefit the State.\nPublic perception of State-sanctioned gaming has also influenced sponsors. Further, a disproportionate amount of the gaming public playing state lotteries is comprised of economically, educationally and/or socially disadvantaged citizens. For example, States that are now dependent upon lottery-generated funds, find themselves in the uncomfortable position of promoting gaming in order to generate revenue with a disproportionate amount of the revenue being raised from the economically, educationally and/or socially disadvantaged. Further, the gaming industry is in a state of flux, with a variety of factors influencing both the financial viability and public perception of State-sanctioned gaming. Significant competitive variables include the proliferation of states in which casinos and/or lotteries are now legal, the increasing reliance of government bodies on the funds generated by state-regulated gaming, the rapid acceptance and use of online and electronic gaming, and the introduction of multi-government-body lotteries with ever-larger jackpots. Thus, State-operated lotteries and games of chance are engaged in a fierce competition for customers with other forms of gaming opportunities.\nGenerally, state lotteries are totally anonymous, with the exception of the recipients of large prizes, who are paid by the State directly. Thus, there is an unfilled need for a system and method for tracking lottery play and proceeds.\nThe following presents a simplified summary of the invention in order to provide a basic understanding of some aspects of the invention. This summary is not an extensive overview of the invention. It is not intended to identify key or critical elements of the invention or to delineate the scope of the invention. Its sole purpose is to present some concepts of the invention in a simplified form as a prelude to the more detailed description that is presented later.\nThe present invention relates to a system and method for tracking lottery play and proceeds. More particularly, the present invention enables sponsors (e.g., state agencies) of gaming opportunities (e.g., lotteries) to track, understand and strategically reward the play of individual players thereby increasing the sponsor\"\"s competitiveness in the market place for gaming funds. Further, the system and method of the present invention increase the social acceptance of the sponsor\"\"s gaming opportunities by allowing a predetermined amount of a registered player\"\"s purchases (e.g., set aside amount) to be deposited in a registered player account for the benefit of the registered player or his/her designate (e.g., personal retirement account, medical savings account, real estate investment account, small business investment account and/or educational savings account).\nAccording to one aspect, of the present invention, a registered player is able to purchase game of chance items (e.g., lottery tickets) from the sponsor through a point of sale computer (e.g., lottery ticket vendor). The registered player identifies himself at the point of sale computer (e.g., using an identification card or driver\"\"s license). The registered player then purchases and receives his game of chance items (e.g., lottery tickets). Information regarding the registered player\"\"s purchases is electronically communicated to the sponsor\"\"s computer system where the registered player\"\"s information is updated to reflect the registered player\"\"s purchases. The registered player\"\"s set aside amount is increased by a predetermined portion of the registered player\"\"s purchase (e.g., ten percent). Once the registered player\"\"s set aside amount reaches a predetermined threshold level (e.g. fifty dollars) in a predetermined period of time (e.g. one year), the registered player\"\"s set aside amount is transferred to an account manager which maintains a registered player account for the benefit of the registered player or his/her designate (e.g., personal retirement account, medical savings account, real estate investment account, small business investment account and/or educational savings account).\nAccording to another aspect of the present invention, a registered player is able to purchase game of chance items (e.g., lottery tickets) directly from the sponsor utilizing a player communication device (e.g., personal computer, telephone or terminal). The registered player identifies himself to the sponsor computer by electronically communicating via the player communication device. The registered player then purchases and receives his game of chance items (e.g., lottery tickets). Information regarding the registered player\"\"s purchases is electronically communicated to the sponsor\"\"s computer system where the registered player\"\"s information can be updated to reflect the registered player\"\"s purchases. The registered player\"\"s set aside amount is increased by a predetermined portion of the registered player\"\"s purchases (e.g., ten percent). Once the registered player\"\"s set aside amount reaches a predetermined threshold level (e.g. fifty dollars) in a predetermined period of time (e.g. one year), the registered player\"\"s set aside amount is transferred to an account manager which maintains a registered player account for the benefit of the registered player or his/her designate (e.g., personal retirement account, medical savings account, real estate investment account, small business investment account and/or educational savings account).\nAccording to yet another aspect of the present invention, a registered player can make direct contributions to his registered player account through a point of sale computer, a player communication device and/or mail. Further, the sponsor is able to utilize information contained in the registered player\"\"s information stored in the sponsor\"\"s computer for targeted marketing campaigns, product development, and/or facilitating registered player income tax reporting (e.g., information regarding registered player\"\"s yearly purchases sent to registered player to provide evidence of a gambling loss deduction and/or enable a potential set off against winnings for a tax year).\nThe present invention also provides methods for tracking game of chance play and proceeds; transferring registered player\"\"s set aside amounts to a registered player account; establishing registered player set asides and registered player accounts."}
{"text": "The difficulties in the development of hair coloring compositions, which can deliver precise long-lasting colors with low hair damage, are in part due to the inherent structure of the hair itself and in part due to the necessary harsh conditions of effective hair coloration processes.\nIn general, the condition and structure of human hair is not regular along the length of the hair shaft. Human hair is subject to various chemical and mechanical treatments such as combing, brushing, shampooing, heating, perming as well as exposure to the sun. As such, the hair at the ends of the hair shaft will generally exhibit signs of damage relative to the new growth close to the scalp. This damage can lead to inconsistent coloration when the hair is dyed due to irregular uptake of the hair-coloring agents along the length of the hair shaft. Thus there is a need for hair coloring compositions which can deliver substantially consistent hair color results throughout the hair.\nOnce the hair has been colored there is a desire for the color to be resistant to the fading action of the sun and other exterior factors and that the color is retained in a consistent manner for a predictable period of time. Thus it would be desirable to develop a hair coloring composition which exhibited reduced fade and provided improved resistance to wash out during regular cleansing regimen. An additional difficulty commonly associated with the dyeing of human hair is the need for dye systems that avoid any adverse effect on the hair and skin of the user.\nTo color human or animal hair using conventional oxidative dye technology it is generally necessary to treat the hair with a mixture of suitable oxidative coloring agents and at least one oxidizing agent and swelling agent at alkaline pHs. Hydrogen peroxide is the most commonly used oxidizing agent. This can lead to undesirable hair damage. These undesirable effects are in part due to the necessary conditions of conventional peroxide treatment, as part of the hair coloring process, which requires high pH ( greater than pH 9), use of swelling agents such as ammonia, extended exposure (from 30 to 60 minutes) and relatively high concentration of oxidizing solutions (up to 30% volume of oxygen) in order to deliver effective hair dyeing and hair melanin bleaching. Thus there is a need for hair coloring compositions, which can oxidize dye precursors and color the hair effectively without damage.\nOxidative dyes and peroxide oxidizing agents can be used to deliver a variety of hair colors to the hair. However such conventional hair coloring compositions comprising oxidative dyes and peroxide oxidizing agent do not deliver against the key consumer needs of color saturation, color vividness, precise initial color predictability, improved wash fastness, improved safety and low hair damage.\nApplicants have now found that the combination of peroxide oxidizing agents with certain transition metal salts or complexes and one or more oxidative coloring agents in hair coloring compositions can deliver excellent initial hair color in a faster time, good wash fastness of the hair color over time, desirable color saturation and vividness attributes, and significantly reduced hair damage. Furthermore, applicants have found that the combination of these dye precursors and transition metal salts or complexes with a source of peroxide at particular levels and ratios in the coloring compositions of the present invention can deliver these excellent hair coloring attributes results with minimal hair damage in a faster time.\nThus it is an object of the present invention to provide improved hair coloring compositions.\nThe following publications relate to the field of the invention:\nU.S. Pat. No. 6,004,355 discloses hair coloring compositions which comprise (a) a water-soluble peroxygen oxidizing agent, (b) an organic peroxyacid precursor oxidizing aid, and (c) one or more oxidative hair coloring agents. The compositions may be used in processes for coloring human or animal hair and may be conveniently provided in kits wherein each of the aforementioned components are individually packaged. This patent, at column 23, lines 20-35, discloses as a ligand, a list which includes 1,4,7-trimethyl-1,4,7-triazacyclononane.\nWO 97/24108, WO 97/24107 and U.S. Pat. No. 6,004,355 disclose a hair bleaching composition comprising: a) a water-soluble peroxygen bleach and (b) a bleaching aid selected from organic peroxyacid bleach precursors and/or preformed organic peroxyacids. The products can provide hair bleaching and in-use efficacy benefits including reduced hair damage at lower pH.\nU.S. Pat. No. 5,873,910 discloses a process for the two-stage oxidation dyeing of keratin fibers by applying to the keratin fibers:\nin a first stage, at least one composition A containing at least one manganese salt and/or a manganese complex,\nin a second stage, at least one composition B having a pH of greater than or equal to 6, and resulting from the extemporaneous mixing of at least one composition B1 containing at least one 4-substituted 1-naphthol and at least one composition B2 containing at least one oxidizing agent, and corresponding multi-compartment dyeing kit.\nThis patent discloses 1,4,7-trimethyl-1,4,7-triazacyclononane as a manganese complex to be used in its process at column 10, lines 20-25.\nDE 19852972 discloses the use of 1,4,7-trimethyl-1,4,7-triazacyclononane in hair dye formation.\nThe present invention is related to a hair coloring composition suitable for the treatment of human or animal hair. According to one aspect of the present invention, there is provided a hair coloring composition comprising a first composition which comprises:\n(a) a dye forming transition metal complex; which is first applied to the hair;\nand a second composition which comprises the following two compositions which are mixed just prior to application to the hair:\n(a) a water-soluble peroxygen oxidizing agent; and\n(b) one or more oxidative hair coloring agents selected from the group consisting of an aromatic diamine, an amino phenol, a polyhydric phenol, a catechol and mixtures thereof.\nAs used herein, % means weight % of the total composition unless otherwise designated. The compositions of the invention are made using known ingredients or with ingredients analogous to those known in the art. The packages or containers to be used with the compositions of the invention are made using known processes and materials or by processes and materials which are analogous to those known in the art.\nCompositions of the invention are used in a two step process for the coloring and bleaching of hair. That is, the first recited composition comprising dye forming transition metal catalysts is contacted with the hair, and then the second composition is made by mixing the following two ingredients just prior to application to the hair; (a) a water-soluble peroxygen oxidizing agent; and (b) one or more oxidative hari coloring agents. After mixture, this composition is then applied to hair. The coloring reaction takes place and the hair is rinsed.\nAs used herein the term xe2x80x9chairxe2x80x9d to be treated may be xe2x80x9clivingxe2x80x9d i.e. on a living body or may be xe2x80x9cnon-livingxe2x80x9d i.e. in a wig, hairpiece or other aggregation of non-living fibres. Mammalian, preferably human hair is preferred. However wool, fur and other melanin containing fibres are suitable substrates for the compositions according to the present invention.\nThe hair coloring compositions can contain, in addition to a mixture of active oxidizing agents, transition metal salts or complexes and oxidative coloring agents such as, by way of example, sequestrants, thickeners, buffers, carriers, surfactants, solvents, antioxidants, polymers non-oxidative dyes and conditioners.\nAs noted above, there is provided a hair coloring composition comprising a first composition which comprises:\n(a) a dye forming transition metal complex or salt;\nwhich is first applied to the hair; and a second composition, which comprises the following two compositions which are mixed just prior to application to the hair:\n(a) a water-soluble peroxygen oxidizing agent; and\n(b) one or more oxidative hair coloring agents selected from the group consisting of an aromatic diamine, an amino phenol, a polyhydric phenol, a catechol and mixtures thereof.\nThe invention also relates to a composition comprising a transition metal wherein the transition metal salt and/or transition metal complex is selected from group 1B, 6B, 7B or 8B metals and an azamacrocycle;\na composition wherein a water-soluble peroxygen oxidizing agent is selected from hydrogen peroxide, sodium perborate, sodium percarbonate, and urea peroxide complexes;\na composition wherein the transition metal complex is formed from a manganese salt and 1,4,7-trimethyl-1,4,7-triazacyclononane;\na composition wherein said manganese complex is bis(octahydro-1,4,7-trimethyl-1H-1,4,7-triazinone-N1, N4, N7)-tri-mxcexc-oxodimanganese bis hexaflurophosphate;\na composition wherein the weight ratio of water-soluble peroxygen oxidizing agent to transition metal complex is in the range of from about 6000:1 to about 1:2;\na composition wherein the weight ratio of water-soluble peroxygen oxidizing agent to transition metal complex is from about 2000:1 to about 1:1;\na composition wherein the weight ratio of water-soluble peroxygen oxidizing agent to transition metal complex is from about from 600:1 to about 60:1;\na composition wherein the said transition metal salt and/or transition metal catalyst are present in a concentration ranging from 0.0001% to 0.2% by weight of metal equivalents relative to the total weight of composition;\na composition wherein the said transition metal salt and/or transition metal catalyst are present in a concentration ranging from 0.001% to 0.1% by weight of metal equivalents relative to the weight of the total composition;\na composition which further comprises a buffering agent;\na composition having one or more oxidative hair coloring agents as described above which has a pH of about 5 to about 11;\na composition as described above which further comprises a surfactant selected from the group consisting of anionic, nonionic, cationic, zwitterionic, amphoteric surfactants and mixtures thereof; and\na hair coloring kit comprising an individually packaged oxidizing component with additional agents, an individually packaged component having one or more hair coloring agents and an individually packaged transition metal salt and/or transition metal complex with additional agents which when mixed forms a composition of the invention.\nAccording to a further aspect of the present invention, there is provided a hair coloring composition comprising:\n(b) from about 0.0001% to about 2% by weight of a transition metal and a second composition which comprises the following two compositions, which are mixed just prior to application to the hair:\n(a) a composition comprising from about 0.01% to about 10% by weight of a water-soluble peroxygen oxidizing agent;\n(b) a composition comprising from about 0.0001% to about 7% by weight of one or more oxidative hair coloring agents;\nThe invention also relates to compositions for coloring human or animal hair wherein the hair coloring composition comprises:\n(a) a composition comprising from about 0.001% to about 1.0% by weight of a transition metal salt and/or transition metal complex;\nwhich is first applied to the hair; and a second composition, which comprises the following two compositions which are mixed just prior to application to the hair:\n(a) a composition comprising from about 0.5% to about 20% by weight of a water-soluble peroxygen bleaching compound\n(b) a composition comprising from about 5% to 50% additional agents selected from surfactants, buffering agents, oxidation bases and/or couplers from about 20% to about 95.5% by weight of an inert diluent\nThe invention also relates to a process for coloring human or animal hair wherein the hair coloring composition comprises applying to the hair the compositions described above.\nIt has been found that compositions of the present invention are unexpectedly effective for lightening and coloring very dark brown or black hair.\nFor example, when (a) a dye forming transition metal complex composition of the present invention is first applied to very dark brown hair or black hair;\nand then a second composition which comprises (a) a water-soluble peroxygen oxidizing agent; and (b) a composition comprising one or more oxidative hair coloring agents the following two compositions (which are mixed just prior to application to the hair) are applied to said very dark hair brown or black hair there is obtained a very dramatic change color of the hair.\nChange in hair color can be measured by means which are known in the art. A customary scale for measuring the change in hair color is defined in terms of the variables: L, a, and b wherein L, a, and b are defined as follows: L indicates the lightness or darkness of the color value. The higher the L, therefore, the lighter the hair, and the more fading that has occurred. When L is 0, the hair is black, and when L is 100, the hair is white. xe2x88x92a and +a represent changes in color tone from green to red. xe2x88x92b and +b represent the changes in color tone from blue to yellow.\nWhat follows now is a description of the ingredients which may be used in the compositions of the invention.\nDye Forming Transition Metal Complex or Salt\nThe compositions of the invention comprise a transition metal containing catalyst for the peroxide oxidizing agents. One suitable type of catalyst is a catalyst system comprising a heavy metal cation of defined catalytic activity, such as copper or manganese cations, an auxiliary metal cation having little or no catalytic activity, such as zinc or aluminum cations, and a sequestrant having defined stability constants for the catalytic and auxiliary metal cations, particularly ethylenediaminetetraacetic acid, ethylenediaminetetra (methylenephosphonic acid) and water-soluble salts thereof. Such catalysts are disclosed in U.S. Pat. No. 4,430,243 which is hereby incorporated by reference.\nOther types of suitable catalysts include the manganese-based complexes disclosed in U.S. Pat. Nos. 5,246,621 and 5,244,594 which is hereby incorporated by reference.\nPreferred examples of these catalysts include MnIV2(xcexcxe2x88x92O)3(1,4,7-trimethyl-1,4,7-triazacyclononane) 2-(PF6)2, MnIII2(xcexcxe2x88x92O)1(xcexcxe2x88x92OAc) 2(1,4,7-trimethyl -1,4,7-triazacyclononane)2-(ClO4)2, MnIIIMnIV4(xcexcxe2x88x92O)1(xcexcxe2x88x92OAc)2-(1,4,7-trimethyl-1,4,7-triazacyclononane) 2-(ClO4)3, and mixtures thereof. Others are described in EP-A-0,549,272 which is hereby incorporated by reference. Other ligands suitable for use herein include 1,5,9-trimethyl-1,5,9-triazacyclododecane, 2-methyl-1,4,7-triazacyclononane, 2-methyl-1,4,7-triazacyclononane, 1,2,4,7-tetramethyl-1,4,7-triazacyclononane, and mixtures thereof.\nFor examples of suitable catalysts see U.S. Pat. Nos. 4,246,612 and 5,227,084 which is hereby incorporated by reference. See also U.S. Pat. No. 5,194,416 which is hereby incorporated by reference which teaches mononuclear manganese (IV) complexes such as Mn(1,4,7-trimethyl-1,4,7-triazacyclononane)(OCH3)3-(PF6). Still another type of suitable catalyst, is disclosed in U.S. Pat. No. 5,114,606, which is hereby incorporated by reference. It is a water-soluble complex of manganese (III) and/or (IV) with a ligand which is a noncarboxylate polyhydroxy compound having at least three consecutive Cxe2x80x94OH groups. Other examples include binuclear Mn complexed with tetra-N-dentate and bi-N-dentate ligands, including N4MnIII(xcexcxe2x88x92O)2MnIVN4)+ and [Bipy2MnIII(xcexcxe2x88x92O)2MnIVbipy2]-(ClO4)3.\nFurther suitable catalysts are described, for example, in EP-A0,408,131 (cobalt complex catalysts), EP-A-0,384,503, and EP-A-0,306,089 which is hereby incorporated by reference (metallo-porphyrin catalysts), and U.S. Pat. No. 4,728,455 all of which are hereby incorporated by reference.\nAlso suitable is (manganese multidentate ligand catalyst), U.S. Pat. No. 4,711,748 and EP-A 0.224,952, (absorbed manganese on aluminosilicate catalyst), U.S. Pat. No. 4,601,845; (aluminosilicate support with manganese and zinc or magnesium salt), U.S. Pat. No. 4,626,373 (manganese ligand catalyst), U.S. Pat. No. 4,119,557 (cobalt chelant catalyst) CA-A-866,191 (transition metal containing salts), U.S. Pat. No. 4,430,243 (chelants with manganese cations and noncatalytic metal cations); and U.S. Pat. No. 4,728,455 (manganese gluconate catalysts); all of which are hereby incorporated by reference.\nOther types of suitable catalysts include the manganese-based salts as described in U.S. Pat. No. 5,873,610 which is hereby incorporated by reference.\nThe compositions of the invention comprise one or more transition metal catalysts. When compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents selected from the group) the concentration of each catalyst is from about 0.0001% to about 1% by weight and is preferably from about 0.001% to about 0.1% by weight.\nThe total combined level of catalysts in the compositions according to the present invention is from about 0.0001% to about 2%, preferably from about 0.001% to about 0.2%, more preferably from about 0.001% to about 0.1% by weight.\nWater-Soluble Peroxygen Oxidizing Agents\nThe compositions of the invention may comprise at least one water-soluble peroxygen oxidizing agent. Water-soluble as defined herein means a peroxygen oxidizing agent compound, which can be substantially solubilized in water.\nThe peroxygen oxidizing agents useful herein are generally inorganic peroxygen materials capable of yielding hydrogen peroxide in an aqueous solution. Water-soluble peroxygen oxidizing compounds are well known in the art and include hydrogen peroxide, inorganic alkali metal peroxides such as sodium periodate and sodium peroxide and organic peroxides such as urea peroxide, melamine peroxide, and inorganic perhydrate salt oxidizing compounds, such as the alkali metal salts of perborates, percarbonates, perphosphates, persilicates, persulphates and the like. These inorganic perhydrate salts may be incorporated as monohydrates, tetrahydrates, and the like. Mixtures of two or more such oxidizing agents can be used if desired. Preferred for use in the compositions according to the present invention is hydrogen peroxide.\nWhen compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents), the peroxygen oxidizing agent is present in the compositions of the present invention at a level of from about 0.01% to about 9%, preferably from about 0.1% to about 6%, more preferably from about 0.2% to about 4% by weight.\nOxidative Hair Coloring Agents\nThe compositions of the present invention may include one or more oxidative hair coloring agents. Such oxidative hair coloring agents are used in combination with the oxidizing systems of the present invention to formulate permanent, hair dye compositions.\nPermanent hair dye compositions as described herein are compositions, which once applied to the hair are substantially resistant to wash-out. Wash-out as described herein is the process by which hair color is removed from the hair over time during normal hair cleansing regimen.\nWhen compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents), the concentration of each oxidative hair coloring agent is from about 0.0001% to about 7% by weight and is preferably from about 0.001% to about 2.0% by weight.\nWhen compared against a combination of all of the compositions of the invention after they have been mixed together (that is, the dye forming transition metal complex; the water-soluble peroxygen oxidizing agent composition; and the composition containing one or more oxidative hair coloring agents), the total combined level of oxidative hair coloring agents in the compositions according to the present invention is from about 0.01% to about 15%, preferably from about 0.01% to about 10%, more preferably from about 0.1% to about 5% by weight.\nOxidative hair coloring agents which can be used in compositions of the invention can be selected from the group consisting of an aromatic diamine, an aminophenol, a polyhydric phenol, a catechol and mixtures thereof. Oxidative hair coloring agents which can also be called oxidative dyes are described in more detail below.\nThe dye forming intermediates used in oxidative dyes can be aromatic diamines, aminophenols and their derivatives. These dye-forming intermediates can be classified as; primary and secondary intermediates. Primary intermediates are chemical compounds, which by themselves will form a dye upon oxidation. The secondary intermediates, are also known as color modifiers or couplers and are used with other intermediates for specific color effects or to stabilize the color.\nThe oxidation dye intermediates which are suitable for use in the compositions and processes herein include aromatic diamines, polyhydric phenols, aminophenols and derivatives of these aromatic compounds (e.g., N-substituted derivatives of the amines, and ethers of the phenols). Primary oxidation dye intermediates are generally colorless molecules prior to oxidation. The oxidation dye color is generated when the primary intermediate is activated and subsequently joined with a secondary intermediate (coupling agent), which is also generally colorless, to form a colored, conjugated molecule.\nIn general terms, oxidation hair dye precursors or intermediates include those monomeric materials which, on oxidation, form oligomers or polymers having extended conjugated systems of electrons in their molecular structure. Because of the new electronic structure, the resultant oligomers and polymers exhibit a shift in their electronic spectra to the visible range and appear colored. For example, oxidation dye precursors capable of forming colored polymers include materials such as aniline, which has a single functional group and which, on oxidation, forms a series of conjugated imines and quinoid dimers, trimers, and the like, ranging in color from green to black. Compounds such as p-phenylenediamine, which has two functional groups, are capable of oxidative polymerization to yield higher molecular weight colored materials having extended conjugated electron systems. Color modifiers (couplers), such as those detailed hereinafter, are preferably used in conjunction with the oxidation dye precursors herein. A representative list of oxidation dye precursors suitable for use is found in Sagarin, xe2x80x9cCosmetic Science and Technology Interscience, Special Edn. Vol 2 pages 308 to 310, which is hereby incorporated by reference.\nThe typical aromatic diamines, polyhydric phenols, aminophenols, and derivatives thereof, described above as primary dye precursors can also have additional substituents on the aromatic ring, e.g. halogen, aldehyde, carboxylic acid, nitro, sulfonic acid and substituted and unsubstituted hydrocarbon groups, as well as additional substituents on the amino nitrogen and on the phenolic oxygen, e.g. substituted and unsubstituted alkyl and aryl groups.\nNonlimiting examples of suitable aromatic diamines, aminophenols, polyhydric phenols and derivatives thereof, respectively, are the following compounds:\no-phenylenediamine,\nm-phenylenediamine,\n2-nitro-p-phenylenediamine,\n1,3,5-triaminobenzene,\n2-hydroxy-p-phenylenediamine,\n2,4-diaminobenzoic acid,\nsodium 2,4-diaminobenzoate,\ncalcium di-(2,4-aminobenzoate),\nammonium 2,4-diaminobenzoate,\ntrimethylammonium 2,4-diaminobenzoate,\ntri-(2-hydroxyethyl)ammonium 2,4-diaminobenzoate,\n2,4-diaminobenzaldehyde carbonate,\n2,4-diaminobenzenesulfonic acid,\npotassium 2,4-diaminobenzenesulfonate,\nN,N-diisopropyl-p-phenylenediamine bicarbonate,\nN,N-dimethyl-p-phenylenediamine,\nN-ethyl-Nxe2x80x2-(2-propenyl)-p-phenylenediamine,\nN-phenyl-p-phenylenediamine,\nN-phenyl-N-benzyl-p-phenylenediamine,\nN-ethyl-Nxe2x80x2-(3-ethylphenyl)-p-phenylenediamine,\n2,4-toluenediamine,\n2-ethyl-p-phenylenediamine,\n2-(2-bromoethyl)-p-phenylenediamine,\n2-phenyl-p-phenylenediamine laurate,\n4-(2,5-diaminophenyl)benzaldehyde,\n2-benzyl-p-phenylenediamine acetate,\n2-(4-nitrobenzyl)-p-phenylenediamine,\n2-(4-methylphenyl)-p-phenylenediamine,\n2-(2,5-diaminophenyl)-5-methylbenzoic acid,\n2-methoxy-p-phenylenediamine,\n2,3-dimethyl-p-phenylenediamine,\n2,5-dimethyl-p-phenylenediamine,\n2-methyl-5-methoxy-p-phenylenediamine,\n2,6-methyl-5-methoxy-p-phenylenediamine,\n3-methyl-4-amino-N,N-diethylaniline,\nN,N-bis-(2-hydroxyethyl)-p-phenylenediamine,\n3-methyl-4-amino-N,N-bis-(2-hydroxyethyl)aniline,\n3-chloro-4-amino-N,N-bis-(2-hydroxyethyl)aniline,\n4-amino-N-ethyl-(piperidonoethyl)aniline,\n3-methyl-4-amino-N-ethyl-xcex2-(piperidonoethyl)aniline,\n4-amino-N-ethyl-N-(morpholinoethyl)aniline,\n4-amino-N-ethyl-N-(acetylaminoethyl)aniline,\n4-amino-N-(methoxyethyl)aniline,\n3-methyl-amino-N-ethyl-N-(2-acetylaminoethyl)aniline,\n4-amino-N-ethyl-N-(mesylaminoethyl)aniline,\n3-methyl-4-amino-N-ethyl-N-(xcex2-mesylaminoethyl)aniline,\n4-amino-N-ethyl-N-(xcex2-sulfoethyl)aniline,\n3-methyl-4-amino-N-ethyl-N-(xcex2-sulfoethyl)aniline,\nN-(4-aminophenyl)morpholine,\nN-(4-aminophenyl)piperidine,\n2,3-dimethyl-p-phenylenediamine,\n2-isopropyl-p-phenylenediamine,\nN,N-bis-(2-hydroxyethyl)-p-phenylenediamine sulfate,\no-aminophenol,\nm-aminophenol,\np-aminophenol,\n2-iodo-p-aminophenol,\n2-nitro-p-aminophenol,\n3,4-dihydroxyaniline,\n3,4-diaminophenol,\n2-hydroxy-4-aminobenzoic acid,\n2-hydroxy-4-aminobenzaldehyde,\n3-amino-4-hydroxybenzenesulfonic acid,\nN,N-diisopropyl-p-aminophenol,\nN-methyl-N-(1-propenyl)-aminophenol,\nN-phenyl-N-benzyl-p-aminophenol sulfate,\nN-methyl-N-(3-ethylphenyl)-p-aminophenol,\n2-nitro-5-ethyl-p-aminophenol,\n2-nitro-5-(2-bromoethyl)-p-aminophenol,\n(2-hydroxy-5-aminophenyl)acetaldehyde,\n2-methyl-p-aminophenol,\n(2-hydroxy-5-aminophenyl)acetic acid,\n3-(2-hydroxy-5-aminophenyl-1-propene,\n3-(2-hydroxy-5-aminophenyl)-2-chloro-1-propene,\n2-phenyl-p-aminophenol palmitate,\n2-(4-nitrophenyl)-p-aminophenol,\n2-benzyl-p-aminophenol,\n2-(4-chlorobenzyl-p-aminophenol perchlorate,\n2-(4-methylphenyl)-p-aminophenol,\n2-(2-amino-4-methylphenyl)-p-aminophenol,\np-methoxyaniline,\ndi-(2-aminoethyl-4-aminophenyl) ether,\ndi-(2-hydroxyethyl-4-aminophenyl) ether,\n(4-aminophenoxy)acetaldehyde,\n(4-aminophenoxy)acetic acid,\n(4-aminophenoxy)methanesulfonic acid,\n1-propenyl-4-aminophenyl ether isobutyrate,\ndi-(2-chloro-1-propenyl-4-aminophenyl) ether,\ndi-(2-nitro-1-propenyl-4-aminophenyl) ether,\ndi-(2-amino-propenyl-4-aminophenyl) ether,\ndi-(2-hydroxy-1-propenyl-4-aminophenyl) ether,\nN-methyl-p-aminophenol,\n3-methyl-4-aminophenol,\n2-chloro-4-aminophenol,\n3-chloro-4-aminophenol,\n2,6-dimethyl-4-aminophenol,\n3,5-dimethyl-4-aminophenol,\n2,3-dimethyl-4-aminophenol,\n2,5-dimethyl-4-aminophenol,\n2-hydroxymethyl-4-aminophenol,\n3-hydroxymethyl-4-aminophenol,\no-hydroxyphenol (catechol),\nm-hydroxyphenol (resorcinol),\np-hydroxyphenol (hydroquinone),\n4-methoxyphenol,\n2-methoxyphenol,\n4-(2-chloroethoxy)phenol,\n4-(2-propenoxy)phenol,\n4-(3-chloro-2-propenoxy)phenol,\n2-chloro-4-hydroxyphenol (2-chlorohydroquinone),\n2-nitro-hydroxyphenol (2-nitrohydroquinone),\n2-amino-4-hydroxyphenol,\n1,2,3-trihydroxybenzene (pyrogallol),\n2,4-dihydroxybenzaldehyde,\n3,4-dihydoxybenzoic acid,\n2,4-dihydroxybenzenesulfonic acid,\n3-ethyl-4-hydroxyphenol,\n3-(2-nitroethyl)-4-hydroxyphenol,\n3-(2-propenyl)-4-hydroxyphenol,\n3-(3-chloro-2-propenyl)-4-hydroxyphenol,\n2-phenyl-4-hydroxyphenol,\n2-(4-chlorophenyl)-4-hydroxyphenol,\n2-benzyl-4-hydroxyphenol,\n2-(2-nitrophenyl)-4-hydroxyphenol,\n2-(2-methylphenyl)-4-hydroxyphenol,\n2-(2-methyl-4-chlorophenyl)-4-hydroxyphenol,\n2-methoxy-4-(1-propenyl)phenol,\n4-hydroxy-3 methoxycinnamic acid,\n2,5-dimethoxyaniline,\n2-methylresorcinol,\naniline,\np-chloroaniline,\np-fluoroaniline,\np-nitroaniline,\np-aminobenzaldehyde,\np-aminobenzoic acid,\nsodium p-aminobenzoate,\nlithium p-aminobenzoate,\ncalcium di-(p-aminobenzoate),\nammonium p-aminobenzoate,\np-aminobenzenesulfonic acid,\npotassium p-aminobenzenesulfonate,\nN-methylaniline,\nN-propyl-N-phenylaniline,\nN-methyl-N-2-propenylaniline,\nN-benzylaniline,\nN-(2-ethylphenyl)aniline,\n4-methylaniline,\n4-(2-bromoethyl)aniline,\n2-(2-nitroethyl)aniline,\n4-aminophenylacetaldehyde,\n4-aminophenylacetic acid,\n4-(2-propenyl)aniline acetate,\n4-(3-bromo-2-propenyl)aniline,\n4-phenylaniline chloroacetate,\n4-(3-chlorophenyl)aniline,\n4-benzylaniline,\n4-(4-iodobenzyl)aniline,\n4-(3-ethylphenyl)aniline,\n4-(2-chloro-ethylphenyl)aniline,\nphenol,\np-chlorophenol,\np-nitrophenol,\np-hydroxybenzaldehyde,\np-hydroxybenzoic acid,\np-hydroxybenzenesulfonic acid,\nethylphenyl ether,\ndi-(2-chloroethylphenyl) ether,\ndi-(2-nitroethylphenyl) ether,\nphenoxyacetaldehyde,\nphenoxyacetic acid,\n3-phenoxy-1-propene,\n3-phenoxy-2-nitro-1-propene,\n3-phenoxy-2-bromo-1-propene,\n4-propyiphenol,\n4-(3-bromopropyl)phenol,\n2-(2-nitroethyl)phenol,\n4-hydroxyphenylacetaldehyde,\n4-hydroxyphenylacetic acid,\n4-(2-propenyl)phenol,\n4-phenyiphenol\n4-benzylphenol,\n4-(3-fluoro-2-propenyl)phenol,\n4-(4-chlorobenzyl)phenol,\n4-(3-ethylphenyl)phenol\n4-(2-chloro-3-ethylphenyl)phenol,\n2,5-xylenol,\n2,5-diaminopyridine,\n2-hydroxy-5-aminopyridine,\n2-amino-3-hydroxypyridine,\ntetraaminopyrimidine,\n1,2,4-trihydroxybenzene\n1,2,4-trihydroxy-5-(C1-C6-alkyl)benzene,\n1,2,3-trihydroxybenzene,\n4-aminoresorcinol,\n1,2-dihydroxybenzene,\n2-amino-1,4-dihydroxybenzene,\n2-amino-4-methoxyphenol,\n2,4-diaminophenol,\n3-methoxy-1,2-dihydroxybenzene,\n4,6-dimethoxy-3-amino-1-hydroxybenzene,\n2,6-dimethyl-4-(p-hydroxyphenyl)amino]phenol\nand salts thereof.\nAdditional oxidation dye couplers suitable for use herein include catechol species and in particular catechol xe2x80x9cdopaxe2x80x9d species which includes dopa itself as well as homologs, analogs and derivatives of DOPA. Other suitable dye precursors are dihydroxyindole (DHI), dihydroxyindolecarboxylic acid (DHICA) and derivatives thereof, indolines and derivatives thereof. Examples of suitable catechol species include cysteinyl dopa, alpha alkyl dopa having 1 to 4, preferably 1 to 2 carbon atoms in the alkyl group, epinephrine and dopa alkyl esters having 1 to 6, preferably 1 to 2 carbon atoms in the alkyl group.\nThe oxidation dye couplers (precursors) can be used herein alone or in combination with other oxidation dye couplers (precursors) mentioned above. The choice of a single dye coupler (precursor) will be determined by the color, shade and intensity of coloration which is desired. The following are preferred oxidation dye couplers (precursors) which can be used herein, singly or in combination, to provide oxidation hair dyes having a variety of shades ranging from ash blonde to black; these are: pyrogallol, resorcinol, p-toluenediamine, o-phenylenediamine, m-phenylenediamine, o-aminophenol, p-aminophenol, 4-amino-2-nitrophenol, nitro-p-phenylenediamine, N-phenyl-p-phenylenediamine, m-aminophenol, 2-amino-3-hydroxytoluene, N,N bis (2-hydroxyethyl)p-phenylenediamine, 4-amino-2-hydroxytoluene, 1,5-dihydroxynapthalene 2,4-diaminoanisole, hydroquinone, 4-amino-2-hydroxytoluene, 2-methyl resorcinol, 2-methyl-5-hydroxyaminophenol, 6-amino-3-hydroxy-toluene, 2,5-diaminotoluene, and 1-phenyl-3-methyl-pyazolone. These can be used in the molecular form or in the form of peroxide-compatible salts, as detailed above.\nApplicant has also found that compositions of the present invention comprising particular oxidizing agents with particular transition metal salts or catalysts with oxidative dyes are valuable for the delivery of good high intensity colors with reduced levels of dye and peroxide. In particular applicant has now found that good hair coloring results can be achieved using the oxidizing systems of the present invention and up to about 84% less dye versus conventional hair coloring compositions. As such the compositions according to the present invention are valuable for the delivery of good initial color and wash fastness over time in addition to having low hair damage attributes.\nSolvents\nWater is the preferred principal diluent for the compositions according to the present invention. As such, the compositions of present invention may also include one or more solvents as additional diluent materials. Generally, the solvent is selected to be miscible with water and innocuous to the skin. Solvents suitable for use herein include C1-C20mono- or polyhydric alcohols and their ethers, glycerine, with monohydric and dihydric alcohols and their ethers preferred. In these compounds, alcoholic residues containing 2 to 10 carbon atoms are preferred. Thus, a particularly preferred group includes ethanol, isopropanol, n-propanol, butanol, propylene glycol, ethylene glycol monoethyl ether, and mixtures thereof.\nThe compositions of the invention may also include the following materials.\nBuffering Agents\nThe coloring compositions of the present invention may have a pH in the range of from about 5 to about 11, more preferably from about 6 to about 8, and especially from about 6 to about 7.\nAs herein before described the preferred coloring compositions of the present invention may contain one or more buffering agents and/or hair swelling agents (HSAs) to adjust the pH to the desired level. Several different pH modifiers can be used to adjust the pH of the final composition or any constituent part thereof.\nThickeners\nThe composition containing one or more oxidative hair coloring agents of the present invention (coloring compositions) may additionally include a thickener at a level of from about 0.05% to about 20%, preferably from about 0.1% to about 10%, more preferably from about 0.5% to about 5% by weight. Thickening agents suitable for use in the compositions of the invention may be selected from the group consisting of oleic acid, cetyl alcohol, oleyl alcohol, sodium chloride, cetearyl alcohol, stearyl alcohol, and synthetic thickeners; and mixtures thereof.\nSurfactant Materials\nThe compositions of the present invention may additionally contain a surfactant system. Suitable surfactants for inclusion in the compositions of the invention generally have a lipophilic chain length of from about 8 to about 22 carbon atoms and can be selected from the group consisting of anionic, cationic, nonionic, amphoteric, zwitterionic surfactants and mixtures thereof.\nMaterials used in compositions of the invention may be employed in the following ranges:"}
{"text": "This invention relates to an anchor device for supporting an element within the vagina of a mammalian female.\nThe present invention is particularly, but not exclusively concerned with the mounting in a suitable location on a mammalian female of a temperature monitoring device which can be used to detect temperature changes indicative of estrus of the female to enable accurate timing of artificial insemination.\nHowever the anchor device with which the present invention is concerned may be used for other purposes such as other telemetry devices or for drug release.\nIt is known that temperature changes can be used to predict the occurance of estrus in mammalian females. In cows a rise in temperature of 0.3 to 0.8 degrees celcius is usually found at estrus. Estrus has a duration of eight to twelve hours and occurs at roughly twenty one day intervals. If the cow does not conceive at first insemination, it is necessary to inseminate again at successive estrus until the cow does conceive. For this reason it is often necessary to continuously monitor the vagina temperature of the cow (to predict the occurance of estrus) for prolonged periods of time. Because vagina temperature varies due to factors other than estrus, it is necessary to compare vagina temperature to the rolling average of temperature for some days previous in order to make a decision as to whether the rise in temperature is significant.\nOther devices are available for detecting estrus and one example comprises a pedometer which detects increases in animal movement which are also indicative of estrus. Other techniques for detection in the past have involved visual observation of the cow behaviour or the use of marker bulls which result in colour marking of the cow upon mounting by the bull during estrus. In some cases the level of progesterone hormone in the milk can be indicative of estrus.. Also some milking machines measure milk temperature and some claims are made that milk temperature can be used to predict the time of estrus. No commercial use of this method has apparently occurred.\nAttempts have been made to manufacture an anchor device which can be located within the vagina of the mammalian female such as a cow. One example is shown in U.S. Pat. No. 4,387,724 (Zartman) and in U.S. Pat. No. 4,377,157 also of Zartman. These devices cannot support a telemetry device of a type comprising a radio transmitter of sufficient dimension to provide a receivable signal over a sufficient distance to enable continual monitoring of the temperature. A further device is known as the Heigro device (Agrophysics, San Francisco, Calif.) which was designed for long term residence in the heifer's vagina. In addition to these devices being unable to receive and support an element of sufficient size for remote telemetry, there is a severe danger that these devices will be expelled and in some examples up to fourty percent were expelled within the first month. It will be appreciated that in order for these devices to be commercially viable, it is necessary that the device remain in place for significant periods of time and also be reusable without the loss or damage which can occur if the device is expelled.\nUp till now, therefore, there has been no anchor device available which enables the relatively large transmitter necessary for continual monitoring to be mounted within the vagina of the mammalian female and which enables the device to remain in position for extended periods of time.\nIn some cases this problem has been overcome by surgical techniques in which the telemetry device is mounted within an incision in the animal. This is of course however disadvantageous in that the surgical invasion is expensive and time consuming and is open to infection.\nAttempts have also been made to detect temperature changes by a measuring device attached to the ear skin of the animal. However it has not been proven that ear skin temperature can be used to effectively detect estrus."}
{"text": "1. Field of Invention\nThis invention relates to a rewinding accessory for monofilament lawn line trimmer devices. More specifically, it is directed to an improved, low cost rewinder that enables users of lawn trimmers to replace a worn monofilament line with maximum efficiency and in the shortest possible time thereby saving time and money.\nMonofilament line used in weed trimmers is available both on spools ready for use or in bulk. The spools are reusable and are characterized generally by a cylindrical or tubular hub with spaced flanges to confine the trimmer line. The size of the spools varies according to the type of trimmer and amount of line on the spool. The line sold in bulk is substantially less expensive than the spooled line but requires winding onto a spool before use. Thus, people who use a great deal of the line may obtain substantial savings by buying the line in bulk and rewinding the line onto a spool. However, the monofilament line used in the trimmers is relatively stiff and resilient making difficult re-spooling by hand. Consequently, a device to hold and turn the spool for rewinding the line onto the spool makes the job of re-spooling simpler and less time consuming. Because the size of spools does vary, the re-spooling device should adapt to a range of spool sizes.\n2. Related Art\nRewinding devices have long been known to the prior art. Illustrative of such rewinders are U.S. Pat. No. 4,717,086 which issued on Jan. 5, 1988 and U.S. Pat. No. 4,762,286 which issued on Aug. 9, 1988. Such rewinders have disc elements adapted to receive an empty spool and a rotating drive crank which facilitates rotation of the spool. Both of these devices uses a spring element to hold the spool in place and compensate for the variety of spool widths. Neither of these devices provides a truly simple remedy for the problem of compensating for various sizes and types of spools. The device shown in U.S. Pat. No. 4,717,086 uses a spring element and a plurality of notches in a shaft in combination with an arm attached to one of the disk elements for adapting to different widths of spools. Likewise, the device shown in U.S. Pat. No. 4,762,286 uses prongs extending from the disc element to engage the spool and facilitate rotation of the spool thereby limiting its universality to various sized spools.\nAnother prior effort to design a spool rewinding device is illustrated by U.S. Pat. No. 4,728,048. This rewinder also uses a spring to hold the spool in place. However, the device uses a disk and a frusto conical shaped element to receive the empty spool and help accommodate for different sized spools. The device uses a power drill to facilitate rewinding.\nOther pertinent examples of rewinders include those revealed in U.S. Pat. Nos. 3,979,833 and 4,164,332, and 4,442,984 none of which deal with lawn trimmer spools.\nThough the above mentioned rewinders may be helpful in rewinding lawn trimmer spools, they can be improved to provide faster rewinding and greater universality with a lower cost, positive drive clamping combination."}
{"text": "1. Field of the Invention\nThe invention relates to a motor control circuit, and more particularly to a single-phase AC motor control circuit for a dryer.\n2. Description of the Related Art\nNowadays, most motors for dryers employ mechanical centrifugal switches with contacts for starting or for switching control. The mechanical centrifugal switch is connected to a starting winding of a motor of a dryer. As the motor of the dryer is started, and rotation speed thereof is 70%-85% of synchronous speed, the mechanical centrifugal switch is switched on under the action of centrifugal force, switches off the starting circuit and operates. During on/off of the mechanical centrifugal switch, arc spark is generated, which affects work life thereof. In addition, the mechanical centrifugal switch must be disposed in the motor of the dryer, which features complex structure, makes it inconvenient to design overall structure of the motor, and increases size of the motor, and material cost of the motor. It is well-known that during a starting phase, the starting winding and the main winding require comparatively large current, if the starting windings, the main windings, and the heated strips of multiple dryers are simultaneously powered up during the starting phase, starting of the motor is to be affected, and a comparatively large workload is formed, which brings about certain impact on a power system, increases starting time, or even makes it impossible for normal starting, and features imperfect control."}
{"text": "A test head is a very important part of a semiconductor IC tester (hereinafter tester) which is connected with the tester by means of thousands of cables to test a semiconductor device (hereinafter DUT) supplied thereto.\nThe test head usually contains several tens of circuit boards (hereinafter board(s)) consisting of ICs and other electronic parts. When the boards loading the ICs and electronic parts are in operation, they generate a great amount of heat. For instance, when the room temperature is 20.degree. C., the temperature inside the housing of the test head rises up to 40.degree. C. even if a built-in fan of the test head operated to cool off inside the housing.\nIn general, the electronic parts such as ICs loaded on the boards are very sensitive to temperature changes in that some functions such as delay time of ICs are easily affected by such changes. Thus, it is desirable that the changes of the room temperature should be minimized.\nMoreover, when the inner temperature of the housing of the test head becomes too high, it causes the high defect rate of the boards, requiring replacement of the boards, adversely affects the reliability of the tester.\nIn addition, the test head is relatively heavy, for instance, it weighs 200 Kg. Thus, it is a very difficult and tiresome task to replace the boards.\nAn example of conventional technology is explained with reference to FIGS. 6 and 7. FIG. 6 is a partially cut-away plan view showing the inside of the test head housing. FIG. 7 is a partially cut-away side view showing the inside of the test head housing.\nInside of the test head housing is arranged to be cooled by a forced air cooling system by a room temperature of a room where the tester is located. A bundle of cables 90 is connected to sockets of boards 80 and 81 in a housing 9 of the test head. The boards 80 and 81 are inserted to board racks 70 and 71. The board rack 70 is positioned between fans 38 and 39. Similarly, the board rack 71 is positioned between fans 40 and 41.\nIn the conventional example, first, the fan 38 generates an air flow by drawing air at the room temperature in the room where the tester is located. Then, the air generated by the fan 38 cools the boards 80 mounted in the board rack 70, wherein the fan 39 blows out the air warmed by the boards 80. In a similar manner, the fan 40 generates an air flow by drawing air. Then, the air generated by the fan 40 cools the boards 81 mounted in the board rack 71, wherein the fan 41 blows out the air.\nIn order to improve the reliability of the test head, the junction temperature of the IC devices mounted on the boards needed to be low.\nIn addition, it is necessary to maintain the inner temperature of the test head and the junction temperature of the IC devices constant and not be affected by the temperature change of the outside atmosphere.\nParticularly, if the temperature inside the test head greatly changes, for instance, more than 3.degree. C., a calibration operation for adjusting delay times throughout the system is required. As a result, a test through-put of the tester is lowered, causing a low test efficiency.\nThus, there is a need in the industry to freely control the temperature in the housing of the test head in a range of, for instance, .+-.1.degree. C.\nThe recent advancement of technology has realized a high density structure wherein a large number of boards are mounted in the test head. However, the advanced technology has also revealed a problem that the temperature inside the test head cannot effectively be controlled by the conventional technology."}
{"text": "Many users post and share information through various social networking services. In an example, a user may use a mobile device to capture a soccer championship game photo of a player scoring a winning goal. The user may upload the soccer championship game photo to an image sharing service for other users to view, comment on, share, and/or the like. The user may tag the soccer championship game photo with a hashtag, such as #championshipgoal, which may facilitate discovery of the championship game photo by other users performing an image search. In another example, multiple users of a microblog discussion service may engage in a microblog discussion regarding a recent political debate. In this way, users may share information through image sharing services, microblog discussion services, social network profiles, and/or a variety of other social networking services."}
{"text": "The present invention relates to a new and distinct cultivar of Guzmania that is a hybrid, hereinafter referred to by the cultivar name xe2x80x98Jazzxe2x80x99. The present invention relates to seeds which are Guzmania cultivar xe2x80x98Jazzxe2x80x99, as well as plants and plant parts produced from these seeds which have all the morphological and physiological characteristics of the Guzmania cultivar xe2x80x98Jazzxe2x80x99. The invention also relates to methods for producing these seeds and plants. Furthermore, the invention relates to progeny plants produced from the cross of xe2x80x98Jazzxe2x80x99 and another Guzmania selection.\nGuzmania is predominantly epiphytic with a few terrestrial species and is native to the tropics. For the most part, species vary in diameter from 7 or 8 inches to 3 or 4 feet and have rosettes of glossy, smooth-edged leaves.\nFloral bracts of Guzmania frequently have brilliant colors and may last for many months. The range of colors for Guzmania is generally from yellow through orange but may also include flame red and red-purple. White or yellow, tubular, three-petalled flowers may also appear on a stem or within the leaf rosette but are usually short-lived.\nGuzmania may be advantageously grown as pot plants for greenhouse or home use. Typically, the plants are shaded from direct sunlight. During the spring to autumn period, the central vase-like part of the leaf rosette is normally filled with water.\nGuzmania is native to tropical America. Leaves of Guzmania are usually formed as basal rosettes which are stiff and entire and in several vertical ranks. Guzmania plants have terminal spikes or panicles which are often bracted with petals united in a tube about as long as the calyx. The ovary is superior and the seeds plumose.\nAsexual propagation of Guzmania is frequently done through the use of tissue culture practices. Propagation can also be from offshoots which can be detached from the mother plant and grown in an appropriate soil or bark mixture.\nMethods for cultivation and crossing of Guzmania are well known. For a detailed discussion, reference is made to the following publications, which are incorporated herein by reference: Benzing, David H., THE BIOLOGY OF THE BROMELIADS, Mad River Press, Inc., Eureka (1980); Zinuner, Karl, BROMELIEN, Verlag Paul Parey, Berlin (1986); and Rauh, Werner, BROMELIEN, Verlag Eugen Ulmer, Stuttgart (1981).\nA Guzinania inbred is produced by brother/sister crossing over several generations to produce a genetically homozygous plant selection. A hybrid cultivar is produced by crossing two genetically distinct inbred lines, collecting seeds produced by the cross, and germinating seeds so-produced to make hybrid plants. The hybrid seeds and plants produced by this method are uniform with respect to their morphological and physiological characteristics.\nA need exists for a greater variety of Guzmania cultivars with attractive ornamental features. Additionally, a need exists for additional Guzrnania hybrid cultivars that can be easily propagated by seed.\nThese and other objectives have been achieved in accordance with the present invention which provides a new cultivar xe2x80x98Jazzxe2x80x99 that is a product of a planned breeding program undertaken by the inventors in Assendelft, The Netherlands, in 1994. The male or pollen parent was a selection of Guzmania lingulata minor identified by Code No. 94206011. The female or seed parent was a selection of Guzmania wittmackii identified by Code No. 94206206.\nBoth parents have a sufficient degree of homozygosity such that the progeny of the cross are genetically and phenotypically uniform. The variety xe2x80x98Jazzxe2x80x99 therefore can be produced by sexual reproduction by crossing 94206011xc3x9794206206 to produce a population of progeny plants, each of which has the combination of characteristics as herein disclosed for the new cultivar. Seeds produced by crossing 94206011xc3x9794206206 have been deposited with the American Type Culture Collection, 10801 University Boulevard, Manassas, Va. 20110-2209, and accorded Deposit Accession No. PTA-3264. 2500 seeds were deposited on Apr. 9, 2001.\nThe cultivar xe2x80x98Jazzxe2x80x99 can also be produced by asexually reproducing progeny from the cross of 94206011xc3x9794206206 because the combination of characteristics as herein disclosed for the new cultivar xe2x80x98Jazzxe2x80x99 are fly fixed and are retained through successive generations of asexual reproduction.\nThis invention relates to seeds which produce Guzmania cultivar xe2x80x98Jazzxe2x80x99.\nThis invention also relates to Guzmania plants, and parts thereof, having all the physiological and morphological characteristics of Guzmania cultivar xe2x80x98Jazzxe2x80x99. This invention relates to a plant produced from seeds which are Guzmania cultivar xe2x80x98Jazzxe2x80x99. This invention also relates to plant parts, such as pollen, seeds or inflorescence produced by Guzmania cultivar xe2x80x98Jazzxe2x80x99.\nThis invention relates to a method of producing seeds which are Guzmania cultivar xe2x80x98Jazzxe2x80x99, by crossing Guzmania wittmackii selection 94206206 as the female parent with Guzmania lingulata minor selection 94206011 as the male parent and the reciprocate cross with 94206206 as the male parent and 94206011 as the female parent and harvesting seeds produced from said crosses.\nThis invention also relates to a method of producing plants having all the physiological and morphological characteristics of the Guzmania cultivar xe2x80x98Jazzxe2x80x99 comprising the steps of (a) crossing Guzmania wittmackii selection 94206206 as the female parent with Guzmania lingulata minor selection 94206011 as the male parent; (b) harvesting seeds produced from said cross; and (c) producing plants from said harvested seeds.\nFurthermore, this invention relates to progeny Guzmania seeds and plants produced from the cross of xe2x80x98Jazzxe2x80x99, as the male or female parent, with another Guzmania selection."}
{"text": "Electrical connector assemblies typically include a first and second electrical connector having connector housings and a plurality of electrical contacts supported by the connector housings. The first electrical connector is configured to be mounted to a first electrical component so as to place the electrical contacts of the first electrical connector in electrical communication with the first electrical component. The second electrical connector is configured to be mounted to a second electrical component so as to place the electrical contacts of the second electrical connector in electrical communication with the second electrical component, The first and second electrical connectors can be mated to each other so as to place the electrical contacts of the first electrical connector in electrical communication with the electrical contacts of the second electrical connector. Thus, when the first and second electrical connectors are mounted to the first and second electrical components, respectively, and mated to each other, the first and second electrical components are placed in electrical communication with each other. Examples of electrical components to which electrical connectors are configured to be mounted include substrates, such as printed circuit boards, and electrical cables."}
{"text": "An important tool in clinical diagnosis of disease is the use of the magnetic resonance imaging (MRI) contrast agents. In MRI, images are acquired by employing radio frequency pulses to excite nuclear spins of a specimen. The observed signal is from the protons of water molecules in the specimen. MRI can capture three-dimensional images without the need for invasive procedures. By imposing one or more orthogonal magnetic field gradients onto a target, an MR image can be obtained using radio frequency pulses to excite nuclear spins as in NMR. Images are based upon the NMR signal from the protons of water molecules, where the signal intensity in a given volume element is a function of the water concentration and relaxation times (T1 and T2).\nMRI has several advantages over other imaging modalities. MRI can image in three dimensions with high spatial and temporal resolution. Unlike while light microscopy and fluorescent microscopy, MRI is not limited by the distance of scattered light onto the cells of interest or dye intensity. It avoids the harmful ionizing radiation of X-ray and CT. Finally, while positron emission tomography (PET) has higher sensitivity, the resolution is much lower than for MRI. MRI can visualize opaque organisms in three dimensions and can follow organisms over time, making it an ideal biological imaging tool.\nSensitivity and intrinsic contrast can be enhanced by using paramagnetic contrast agents, such as the commonly used paramagnetic Gd(III) ion, that decreases the local T1 relaxation of nearby water protons (Caravan et al. Chem Rev 1999, 99, 2293-352., herein incorporated by reference in its entirety). Images derived from changes in T1 regions that are associated with a Gd(III) ion have a higher signal intensity (Aime et al. Curr Pharm Biotechnol 2004, 5, 509-18., herein incorporated by reference in its entirety). The ability of a contrast agent to decrease T1 and therefore increase signal intensity at a given concentration is relaxivity (mM−1 s−1). High relaxivity agents result in area of increased signal.\nFree Gd(III) ions are toxic to biological systems, and a suitable ligand or chelate must bind the lanthanide to form a nontoxic complex. For many years, Gd(III) contrast agents have been used in a chelated form to eliminate toxicity in humans. Several factors influence the stability of chelate complexes including enthalpy and entropy effects (e.g., number, charge and basicity of coordinating groups, ligand field, and conformational effects). Recently, there has been concern associated with the use of Gd(III) agents due to an apparent link to a disabling condition called NSF (Kurtkoti & Hiremagalur. Nephrology 2008, 13, 235-41., Kay Ann Rheum Dis 2008, 67 Suppl 3, iii66-9., herein incorporated by reference in their entirety). Many patients with this condition experience a thickening of the skin that inhibits joint movement. Clinical data on known cases of NSF have revealed that the condition is only present in patients with renal failure (low pH) and only when certain classes of contrast agent are used. Three clinically approved contrast agents have been associated with the onset of NSF: OMNISCAN, MAGNEVIST and OPTIMARK. These are all linear Gd(III) chelates based on the structure of DTPA.\nMacrocyclic chelates have higher thermodynamic stability constants and have not been associated with NSF. The reduced thermodynamic stability constants (and the presence of an amide) in linear chelates is thought to be the cause of NSF when Gd(III) is released from the chelate and displaced by other naturally occurring metals. The medical community has studied the risks of using gadolinium in patients with renal failure and has published guidelines to minimize the risk (Shellock & Spinazzi, AJR Am J Roentgenol 2008, 191, 1129-39., herein incorporated by reference in its entirety).\nGd(III) ions are toxic to living tissues, presumably due to binding to calcium channels and therefore, it must be chelated to reduce the bioavailability. These chelates are synthetically versatile and provide the means to attach targeting moieties (Allen & Meade. Met Ions Biol Syst 2004, 42, 1-38., herein incorporated by reference in its entirety).\nMammary epithelial cells express the progesterone receptor (PR) and estrogen receptors (ER) (Ismail et al. Steroids 2003, 68, 779-87., herein incorporated by reference in its entirety). The PR is present in two distinct isoforms both derived from the same gene, PRA and PRB. Each subtype is critical to mammary gland lobuloalveolar development and epithelial differentiation (Lanari & Molinolo Breast Cancer Res 2002, 4, 240-3., herein incorporated by reference in its entirety). The receptor consists of several regions that serve as functional units such as the DNA binding domain (DBD), the ligand binding domain (LBD), and transcriptional activation domains (AFs). The expression of these receptors is a critical parameter typically examined using immunohistochemistry in biopsies of human breast cancers (Jacobsen et al. J Mammary Gland Biol Neoplasia 2003, 8, 257-68., herein incorporated by reference in its entirety).\nThe presence of both receptors correlates significantly with the survival rate of breast cancer patients (Hopp et al. Clin Cancer Res 2004, 10, 2751-60., herein incorporated by reference in its entirety). The PR is an estrogen-regulated gene that becomes activated and expressed in the presence of estradiol and ER. Therefore, it is not surprising that treatment with tamoxifen (an anti-estrogenic therapy) reduces PR. Decreased PR correlates with tamoxifen resistance, although the mechanism of resistance is still being debated (Arpino et al. J Natl Cancer Inst 2005, 97, 1254-61., herein incorporated by reference in its entirety). Tumors that are ER+/PR− are considered more metastatic and aggressive than PR+ tumors and correlate with a lower survival rate (Cui et al. J Clin Oncol 2005, 23, 7721-35., herein incorporated by reference in its entirety).\nAn important prognostic marker is the presence of ER+/PR− tumors because these cancers respond much better to aromatase inhibitors than ER+/PR+ tumors that can be effectively treated with tamoxifen (Fuqua et al. J Clin Oncol 2005, 23, 931-2; author reply 932-3., Osborne et al. Breast 2005, 14, 458-65., herein incorporated by reference in their entireties). In addition, expression of PR may also reflect activation of the growth factor pathway Her2/neu. Since monoclonal antibodies directed against the Her2/neu receptor are being developed for breast cancer, knowing the PR status might help determine if a patient will respond to these treatments (Montemurro & Aglietta, Clin Breast Cancer 2005, 6, 77-80., herein incorporated by reference in its entirety). Therefore, non-invasively delineating whether or not a mammary cancer expresses the PR may be crucial to determining the best chemotherapeutic agent for the patient and ultimately improve survival.\nThe progression of endometrial cancer resembles that of breast cancer in regards to the expression of progesterone receptors. For example, the loss of both PRA and PRB is associated with a poor prognosis and inversely correlated with disease free survival (Boruban et al. Eur J Cancer Prey 2008, 17, 133-8., Uharcek. Obstet Gynaecol Res 2008, 34, 776-83., herein incorporated by reference in their entireties). When each isoforms of the receptor is analyzed separately, the correlation to disease progression is less clear (Arnett-Mansfield et al. Cancer Res 2001, 61, 4576-82., herein incorporated by reference in its entirety); however, the contrast agent would bind to all available progesterone receptors and therefore the overall loss of PR is more critical for imaging purposes. In cell lines, a reduction in the level of PRs is associated with increases in genes that regulate invasion (Miyamoto et al. J Steroid Biochem Mol Biol 2004, 92, 111-8., Saito et al. Cancer Sci 2006, 97, 1308-14., herein incorporated by reference in their entireties). Further, in tissue samples PR is inversely correlated to Ki67, a marker of cellular proliferation. Progesterone is also a therapeutic agent for endometrial cancers with many patients receiving progesterone to slow the growth of their cancer (Ito et al. Endocr J 2007, 54, 667-79., herein incorporated by reference in its entirety).\nRecent reports have addressed the use of MR imaging for analyzing breast tumors (Lehman et al. N Engl J Med 2007, 356, 1295-303., Tozaki. Breast Cancer 2008, 15, 205-11., Lee et al. Radiology 2008, 246, 763-71., herein incorporated by reference in their entireties). Although the majority of physicians conclude that traditional mammography is the best for routine screening of the general population, MR imaging is increasingly used in tumor imaging. For patients with familial risk of breast cancer lesions tend to form quickly and have varying appearance using mammography. When a patient has a positive mammography and biopsy, MR imaging is a second line technique to discover other lesions and identify lesions in the contralateral breast. MR can be helpful for guiding biopsies so that the needle is inserted directly into the cancerous area for accurate results.\nMR imaging is valuable for determining if a patient is responding to therapy. Response to therapy is one of the critical areas that a targeted steroid-based contrast agent can be used because many drugs for breast cancers down regulate estrogen inducible genes, such as the progesterone receptor. For uterine cancers in particular, progesterone is often given to the patient as part of treatment and in the case of the PR-imaging agent, the technology would possibly be both therapeutic and diagnostic (theranostic).\nProgesterone agents have been developed for positron emission topography (PET) imaging with success in targeting breast cancer cells and tissues in rat models (Zhou et al. J Med Chem 2006, 49, 4737-44., Vijaykumar et al. A. J Org Chem 2002, 67, 4904-10., Pomper et al. J Med Chem 1988, 31, 1360-3., herein incorporated by reference in their entireties). Metabolic conversion of reported progestin based PET agents prevented the application of these probes in humans (Dehdashti et al. J Nucl Med 1991, 32, 1532-7., herein incorporated by reference in its entirety).\nIn vivo imaging agents could provide a tool for basic scientific investigations into the etiology of disease by providing size and molecular profiles of tumors without the need to euthanize the animal. Mouse models of cancer and uterine tumors are an essential component of understanding how to prevent and treat disease. Many models could be improved by applying imaging techniques such that tumors could develop and differentiate without the need to remove the tumor mass directly.\nFor example, models of ductal carcinoma in situ (DCIS) are available, but understanding when the lesion forms, where, and whether it is steroid responsive is difficult to accomplish without sacrificing the animal and performing a dissection. However, if contrast agents were applied, small lesions could be identified very early, allowing one to determine if they form invasive cancers, and then categorizing the cancer as steroid responsive or unresponsive. Progesterone receptor positive breast and uterine cancer cells can be subcutaneously injected into nude mice and produce solid tumors monitored with magnetic resonance imaging (Zong et al. Magn Reson Med 2005, 53, 835-42., Preda et al. J Magn Reson Imaging 2004, 20, 865-73., herein incorporated by reference in their entireties). The breast cancer cell lines typically used for such tumors include T47D and MCF7 cell lines, both of which express estrogen receptor (ER) and progesterone receptor (PR) (Hoffmann et al. J Natl Cancer Inst 2004, 96, 210-8., herein incorporated by reference in its entirety). For uterine cancers, the Ishikawa cell line lacks receptors and stable clones of the cell line with the PR gene integrated allow the investigator to analyze both receptor positive and negative tumors. Xenografted tumors visibly protrude from the mouse but may be analyzed earlier and with more ease using MR imaging (Bhujwalla et al. Neoplasia 2001, 3, 143-53., herein incorporated by reference in its entirety). These nude mouse tumor models will be used to study the targeting ability of progesterone based contrast agents to image PR+ receptor positive tumors and for quantitative imaging to determine tumor response to drug therapies."}
{"text": "Three-dimensional vertical NAND strings having one bit per cell are disclosed in an article by T. Endoh et al., titled “Novel Ultra High Density Memory With A Stacked-Surrounding Gate Transistor (S-SGT) Structured Cell”, IEDM Proc. (2001) 33-36."}
{"text": "Programmable logic devices (“PLDs”) are a well-known type of integrated circuit that can be programmed to perform specified logic functions. One type of PLD, the field programmable gate array (“FPGA”), typically includes an array of programmable tiles. These programmable tiles can include, for example, input/output blocks (“IOBs”), configurable logic blocks (“CLBs”), dedicated random access memory blocks (“BRAMs”), multipliers, digital signal processing blocks (“DSPs”), processors, clock managers, delay lock loops (“DLLs”), and so forth. As used herein, “include” and “including” mean including without limitation.\nEach programmable tile typically includes both programmable interconnect and programmable logic. The programmable interconnect typically includes a large number of interconnect lines of varying lengths interconnected by programmable interconnect points (“PIPs”). The programmable logic implements the logic of a user design using programmable elements that can include, for example, function generators, registers, arithmetic logic, and so forth.\nThe programmable interconnect and programmable logic are typically programmed by loading a stream of configuration data into internal configuration memory cells that define how the programmable elements are configured. The configuration data can be read from memory (e.g., from an external PROM) or written into the FPGA by an external device. The collective states of the individual memory cells then determine the function of the FPGA.\nAnother type of PLD is the Complex Programmable Logic Device, or CPLD. A CPLD includes two or more “function blocks” connected together and to input/output (“I/O”) resources by an interconnect switch matrix. Each function block of the CPLD includes a two-level AND/OR structure similar to those used in Programmable Logic Arrays (“PLAs”) and Programmable Array Logic (“PAL”) devices. In CPLDs, configuration data is typically stored on-chip in non-volatile memory. In some CPLDs, configuration data is stored on-chip in non-volatile memory, then downloaded to volatile memory as part of an initial configuration (programming) sequence.\nFor all of these programmable logic devices (“PLDs”), the functionality of the device is controlled by data bits provided to the device for that purpose. The data bits can be stored in volatile memory (e.g., static memory cells, as in FPGAs and some CPLDs), in non-volatile memory (e.g., FLASH memory, as in some CPLDs), or in any other type of memory cell.\nOther PLDs are programmed by applying a processing layer, such as a metal layer, that programmably interconnects the various elements on the device. These PLDs are known as mask programmable devices. PLDs can also be implemented in other ways, e.g., using fuse or antifuse technology. The terms “PLD” and “programmable logic device” include but are not limited to these exemplary devices, as well as encompassing devices that are only partially programmable. For example, one type of PLD includes a combination of hard-coded transistor logic and a programmable switch fabric that programmably interconnects the hard-coded transistor logic.\nVoltage-controlled oscillators (“VCOs”) are building blocks in many high performance analog circuits and systems. For example, VCOs are used in frequency synthesizers, clock multiplication units (“CMUs”), clock and data recovery (“CDR”) circuits, and phase-locked-loops (“PLLs”), among other circuits. VCOs have been implemented with CMOS. The benefits of leveraging high-volume CMOS fabrication for analog and RF devices would be useful. For example, it would be useful if such integration of VCOs was complete, namely forming a VCO entirely in a single-chip, such as in an SOC implementation for example. Some design considerations for VCOs are tuning range and jitter performance. In PLDs, these design considerations may be more stringent because a single PLD may be used to support multiple protocols and applications with different clock frequency performance parameters and with stringent jitter performance parameters.\nAccordingly, it would be desirable and useful to provide a high gain varactor circuit capable of meeting stringent jitter performance parameters over a relatively large tuning range to provide a VCO suitable for an application where a single PLD supports multiple protocols and applications.\nA VCO “tank” circuit is a well-known circuit for generating oscillations. VCO circuits, like the tank circuit, are known as LC VCOs, where the L and C respectively indicate inductor and capacitor. A tank capacitance “C” in an LC VCO is a combination of a fixed capacitance and a tunable capacitance. The tunable capacitance conventionally includes a varactor. However, a conventional varactor has a relatively limited capacitive range. Furthermore, as is known, CMOS processes may vary from time to time and thus a capacitance-voltage (“C-V”) curve may shift. Heretofore, such a process shift may reduce the useable range of the varactor capacitance.\nAccordingly, it would be both useful and desirable to provide a varactor circuit that is programmable for a wide range of capacitances. Moreover, it would be useful and desirable to enhance a varactor's C-V curve gain beyond a device intrinsic C-V curve property. It would additionally be useful and desirable if such varactor circuit could be adjusted to offset process variation."}
{"text": "It is known that various natural alumino-silicate minerals are beneficial to animals, and even humans, for the treatment of a range of digestive aliments and to improve the efficiency of animal feeds. Kaolinite is widely used for the treatment of acid stomach and various digestive imbalances. Bentonite is also used for similar treatment in horses and cattle. More recently, natural and manufactured zeolites have been tested and found to have some varying degree of effectiveness in improved feed conversion, treating digestive ailments, improved animal condition and improved healing of injuries. However, the effectiveness of natural zeolites has in general been found to be rather poorer than the effectiveness of manufactured zeolites.\nNatural zeolites are alumino-silicate minerals with various forms of an open box-work crystal structure. There are some 41 known natural zeolite minerals of which clinoptilolite; mordenite, erionite, chabazite and stellerite are but a small number. Molecular bonds within these crystal structures create an attractive force that pulls cations, anions and other small charged particles into the crystal structure. Some powdery natural zeolites are treated by a sintering process to improve their handling characteristics and some are dosed with metal compounds, such as manganese, to alter their reaction to various processes. In the inventor's experience, most crystalline zeolites are denigrated by excessive heat treatment. It was considered by the inventor that by using some form of irradiation it may be possible to enhance the internal molecular bonds and thereby make zeolite minerals and other natural alumino-silicate minerals more attractive to small charged particles.\nIt was also considered by the inventor that improving the reactivity of natural zeolites and other natural alumino-silicate minerals in this way might produce enhanced animal husbandry, veterinary and pharmaceutical products, as well as improved methods of treatment of both human and non human animals, and particularly mammals.\nIt was also considered by the inventor that subjecting natural alumino-silicate minerals to the type of radiation described may disturb the molecular structure of the alumino-silicate minerals and render some of the silicon atoms bio-available."}
{"text": "Newly synthesized eukaryotic mRNA molecules, known as primary transcripts or pre-mRNA, made in the nucleus, are processed before or during transport to the cytoplasm for translation. A methylated cap structure, consisting of a terminal nucleotide, 7-methylguanylate, is added to the 5' end of the mRNA in a 5'--5' linkage with the first nucleotide of the mRNA sequence.\nAn approximately 200-250-base sequence of adenylate residues, referred to as poly(A), is added posttranscriptionally to a site that will become the 3' terminus of the mRNA, before entry of the mRNA into the cytoplasm. This is a multistep process which involves assembly of a processing complex, then site-specific endonucleolytic cleavage of the precursor transcript, and addition of a poly(A) \"tail.\" In most mRNAs the polyadenylation signal sequence is a hexamer, AAUAAA, located 10 to 30 nucleotides in the 5' direction (upstream) from the site of cleavage (5'-CA-3') in combination with a U or G-U rich element 3' to the cleavage site. Multiple poly(A) sites may be present on a given transcript, of which only one is used per transcript, but more than one species of mature mRNA transcript can be produced from a given pre-mRNA via use of different poly(A) sites. It has recently been shown that stable mRNA secondary structure can affect the site of polyadenylation of an RNA construct in transfected cells. Klasens et al., Nuc. Acids Res., 1998, 26, 1870-1876. It has also been found that which of multiple polyadenylation sites is used can affect transcript stability. Chu et al., J. Immunol., 1994, 153, 4179-4189. Antisense modulation of mRNA polyadenylation has not previously been reported.\nThe next step in mRNA processing is splicing of the mRNA, which occurs in the maturation of 90-95% of mammalian mRNAs. Introns (or intervening sequences) are regions of a primary transcript (or the DNA encoding it) that are not included in a finished mRNA. Exons are regions of a primary transcript that remain in the mature mRNA when it reaches the cytoplasm. The exons are \"spliced\" together to form the mature mRNA sequence. Intron-exon junctions are also referred to as \"splice sites\" with the 5' side of the junction often called the \"5' splice site,\" or \"splice donor site\" and the 3' side the \"3' splice site\" or \"splice acceptor site.\" \"Cryptic\" splice sites are those which are less often used but may be used when the \"usual\" splice site is blocked or unavailable.\" Alternative splicing, i.e., the use of various combinations of exons, often results in multiple mRNA transcripts from a single gene.\nA final step in RNA processing is turnover or degradation of the mRNA. Differential mRNA stabilization is one of several factors in the rate of synthesis of any protein. mRNA degradation rates seem to be related to presence or absence of poly(A) tails and also to the presence of certain sequences in the 3' end of the mRNA. For example, many mRNAs with short half-lives contain several A(U).sub.n A sequences in their 3'-untranslated regions. When a series of AUUUA sequences was inserted into a gene not normally containing them, the half life of the resulting mRNA decreased by 80%. Shaw and Kamen, Cell, 1986, 46, 659. This may be related to an increase of nucleolytic attack in sequences containing these A(U).sub.n A sequences. Other mediators of mRNA stability are also known, such as hormones, translation products (autoregulation/feedback), and low-molecular weight ligands.\nAntisense compounds have generally been used to interfere with protein expression, either by interfering directly with translation of the target molecule or, more often, by RNAse-H-mediated degradation of the target mRNA. Antisense interference with 5' capping of mRNA and prevention of translation factor binding to the mRNA by oligonucleotide masking of the 5' cap have been disclosed by Baker et al. (WO 91/17755).\nAntisense oligonucleotides have been used to modulate splicing, particularly aberrant splicing or splicing of mutant transcripts, often in cell-free reporter systems. A luciferase reporter plasmid system has been used to test the ability of antisense oligonucleotides targeted to the 5' splice site, 3' splice site or branchpoint to inhibit splicing of mutated or wild-type adenovirus pre-mRNA sequences in a reporter plasmid. Phosphorothioate oligodeoxynucleotides that can support RNAse H cleavage were found to be better inhibitors of expression of the wild-type adenovirus construct than the 2'-methoxy phosphorothioates that cannot support RNase H, although the reverse was true for oligonucleotides targeted to an adenovirus construct containing human .beta.-globin splice site sequences. Hodges and Crooke, Mol. Pharmacol., 1995, 48, 905-918.\nAntisense oligonucleotides have been used to target mutations that lead to aberrant splicing in several genetic diseases. Use of antisense compounds to correct aberrant processing of mutated mRNA sequences is not comprehended by the present invention. Altering, i.e., controlling, the behavior of a cell, particularly the response of a cell to a stimulus, by antisense modulation of \"wild-type\" or native mRNA processing, the subject of the present invention, has not been described previously.\nPhosphorothioate 2'-O-methyl oligoribonucleotides, have been used to target the aberrant 5' splice site of the mutant .beta.-globin gene found in patients with .beta.-thalassemia, a genetic blood disorder. Aberrant splicing of mutant .beta.-globin mRNA was blocked in vitro in vector constructs containing thalassemic human .beta.-globin pre-mRNAs using 2'-O-methyl-ribo-oligonucleotides targeted to the branch point sequence in the first intron of the mutant human .beta.-globin pre mRNAs. 2'-O-methyl oligonucleotides are used because they are stable to RNAses and form stable hybrids with RNA that are not degraded by RNAse H. Dominski and Kole, Proc. Natl. Acad. Sci. USA, 1993, 90, 8673-8677. A review article by Kole discusses use of antisense oligonucleotides targeted to aberrant splice sites created by genetic mutations such as .beta.-thalassemia or cystic fibrosis. It was hypothesized that blocking a splice site with an antisense oligonucleotide will have similar effect to mutation of the splice site, i.e., redirection of splicing. Kole, Acta Biochimica Polonica, 1997, 44, 231-238. Oligonucleotides targeted to the aberrant .beta.-globin splice site suppressed aberrant splicing and at least partially restored correct splicing in HeLa cells expressing the mutant transcript. Sierakowska et al., Nucleosides & Nucleotides, 1997, 16,1173-1182; Sierakowska et al., Proc. Natl. Acad. Sci. USA, 1996, 93, 12840-44. U.S. Pat. No. 5,627,274 discloses and WO 94/26887 discloses and claims compositions and methods for combating aberrant splicing in a pre-mRNA molecule containing a mutation, using antisense oligonucleotides which do not activate RNAse H.\nModulation of mutant dystrophin splicing with 2'-O-methyl oligoribonucleotides has been reported both in vitro and in vivo. In dystrophin Kobe, a 52-base pair deletion mutation causes exon 19 to be skipped during splicing. An in vitro minigene splicing system was used to show that a 31-mer 2'-O-methyl oligoribonucleotide complementary to the 5' half of the deleted sequence in dystrophin Kobe exon 19 inhibited splicing of wild-type pre-mRNA. Takeshima et al., J. Clin. Invest., 1995, 95, 515-520. The same oligonucleotide was used to induce exon skipping from the native dystrophin gene transcript in human cultured lymphoblastoid cells.\nDunckley et al., (Nucleosides & Nucleotides, 1997, 16, 1665-1668) describes in vitro constructs for analysis of splicing around exon 23 of mutated dystrophin in the mdx mouse mutant, a model for Duchenne muscular dystrophy. Plans to analyze these constructs in vitro using 2' modified oligos targeted to splice sites within and adjacent to mouse dystrophin exon 23 are discussed, though no target sites or sequences are given.\n2'-O-methyl oligoribonucleotides were subsequently used to correct dystrophin deficiency in myoblasts from the mdx mouse. An antisense oligonucleotide targeted to the 3' splice site of murine dystrophin intron 22 caused skipping of the mutant exon and created a novel in-frame dystrophin transcript with a novel internal deletion. This mutated dystrophin was expressed in 1-2% of antisense treated mdx myotubes. Use of other oligonucleotide modifications such as 2'-O-methoxyethyl phosphodiesters are disclosed. Dunckley et al. (Human Mol. Genetics, 1998, 5, 1083-90).\nPhosphorothioate oligodeoxynucleotides have been used to selectively suppress the expression of a mutant .alpha.2(I) collagen allele in fibroblasts from a patient with osteogenesis imperfecta, in which a point mutation in the splice donor site produces mRNA with exon 16 deleted. The oligonucleotides were targeted either to the point mutation in the pre-mRNA or to the defectively spliced transcript. In both cases mutant mRNA was decreased by half but the normal transcript is also decreased by 20%. This was concluded to be fully accounted for by an RNAse H-dependent mechanism. Wang and Marini, J. Clin Invest., 1996, 97, 448-454.\nA microinjection assay was used to test the antisense effects on SV40 large T antigen (TAg) expression of oligonucleotides containing C-5 propynylpyrimidines, either as 2'-O-allyl phosphodiester oligonucleotides, which do not elicit RNAse H cleavage of the target, or as 2'-deoxy phosphorothioates, which do elicit RNAse H cleavage. Oligonucleotides targeted to the 5' untranslated region, translation initiation site, 5' splice junction or polyadenylation signal of the TAg transcript were injected into the nucleus or cytoplasm of cultured cells. The only 2'-O-allyl (non-RNAse H) oligonucleotides which were effective at inhibiting T-antigen were those targeted to the 5' untranslated region and the 5' splice junction. The 2'-O-allyl phosphodiester/C-5 propynylpyrimidine oligonucleotides, which do not elicit RNAse H, were 20 fold less potent than the oligodeoxynucleotides which had the ability to recruit RNAse H. The authors concluded that the duplexes formed between the RNA target and the 2'-O-allyl phosphodiester/C-5 propynylpyrimidine oligonucleotides dissociate rapidly in cells. Moulds et al., 1995, Biochem., 34, 5044-53. Biotinylated 2'-O-allyloligoribonucleotides incorporating 2-aminoadenine bases were targeted to the U2 small nuclear RNA (snRNA), a component of the spliceosome, in HeLa nuclear extracts. These inhibited mRNA production with a concomitant accumulation of splicing intermediates. The present invention is directed to antisense compounds targeted to mRNA.\nUse of antisense compounds to block or regulate mRNA polyadenylation has not previously been described. Regulation of mRNA stability using antisense oligonucleotides targeted to RNA sequences involved in RNA turnover or degradation has also not been previously described.\nThere is, therefore a continued need for compositions and methods for altering the behavior of a cell, particularly the response of a cell to a stimulus, by modulation of normal mRNA processing. The present invention provides antisense compounds for such modulation. The compositions and methods of the invention can be used in therapeutics, including prophylaxis, and as research tools."}
{"text": "During normal driving, including some acceleration, fuel is supplied along with air to the manifold of the engine. At all such times, some of the fuel wets the walls of the intake manifold and a quantity of fuel is usually present in a liquid form. This condition exists even though the pressure in the manifold is substantially below atmospheric e.g. a vacuum of 15 to 18 inches mercury at cruising conditions.\nWhen a vehicle driver removes his foot from the gas pedal, the throttle closes to the curb idle position and the manifold vacuum becomes much greater such as, for example, 25 inches of mercury. This sharply reduced absolute pressure causes at least a portion of a liquid fuel in the intake manifold to flash to the vapor stage. Since air flow has been sharply reduced to a curb idle quantity, and since idle fuel is still being supplied by the idle fuel passage of the carburetor, the \"vaporized\" fuel results in an excess amount of fuel being supplied the engine. Such an excess can be detrimental in the sense that emissions are increased during such a transient condition. Such increased emissions are primarily unburned hydrocarbons (HC) and carbon monoxide (CO). Oxides of nitrogen (NOX) probably are not affected.\nSome emission control devices being considered for future automotive use require a constant air/fuel ratio delivered to the engine under all driving conditions. Rapid increases in intake manifold vacuum result in a temporarily rich air/fuel mixture due to the evaporation of fuel which has previously condensed on the walls of the manifold. It is desirable to lean this temporarily rich condition."}
{"text": "The Internet is a popular means for distributing and viewing video content from a wide variety of media sources, as well as user generated video content. A variety of different types of videos from various sources may be displayed on a computerized device. For instance, displayed videos may include video advertisements, recorded broadcast videos (e.g., news broadcasts, cable and/or television shows and stations, etc.), user generated videos (e.g., videos created and uploaded to a particular website by an individual), and even real-time data streaming videos (e.g., live news coverage, video conferencing, etc.).\nA user will often open and view a video within a window of a computerized device. The window may contain one or more programs that allow executing several programs simultaneously (i.e., allows multitasking), and may include one or more tabs within each window. After the video is opened in the first tab, and perhaps even while the video is being played, the user may open another tab within such window, or even open a different window on the computerized device. For instance, the user may receive a chat message and go to chat, or decide to open another tab, or even go to email. In such an event, the video that is opened and being played in the first tab becomes occluded (i.e., obstructed), such that, the user is not able to view such video. The occlusion of the video may also cause the video to stop playing all together.\nOnce a video becomes occluded the user will miss out watching the content of the Video until the user goes back to the window that has the video. By the time the user goes back to this window/tab, the video would have advanced much further and so the user will have to replay the video from where he had left earlier, typically through a manual process. In instances where the video was buffered from the Internet, the user will have to buffer that piece of the video again.\nIn other instance, once a video becomes occluded the user may need to browse back and forth amongst the various open tabs and/or windows to view the opened video. When the tab or window having the previously viewed video has been reopened, the user may find that the video has stopped playing, or the user may need to search for the video if the page has moved onto a next video within a sequence of videos. Often, if the user needs to reopen the video, such video does not open at the location where the video stopped playing. Rather, the video is loaded from its beginning and the user must watch the entire video again.\nIn other instances, the user may open several windows on the computerized device and have different videos playing simultaneously in the different windows. The user may also need to browse back and forth amongst the different open windows to watch the videos being played thereon, such that, the user will miss portions of videos being played on the windows not being currently viewed.\nAccordingly, there continues to be a need in the art for computer-implemented methods, apparatus, and computer program codes for rendering videos watched in a multi-window computer system where people watch items over multiple windows and/or tabs."}
{"text": "Polyacetal is a well known engineering resin which exhibits excellent chemical resistance properties, as well as balanced mechanical properties. Polyacetal resin has therefore been used in a variety of end-use applications. However, since polyacetal resins also tend to exhibit high coefficients of shrinkage and linear expansion during molding due to the high crystallinity of the resin, it is known to have fairly poor dimensional accuracy. In addition, the weather resistance and acid resistance properties of polyacetal resin are known to need improvements.\nThese disadvantageous properties of polyacetal thus make it difficult to use the resin in certain end-use applications, such as the material from which exterior components of automobiles and the like (which are exposed directly to atmospheric conditions) may be made. Although some attempts have been made to improve the weather resistance properties of polyacetal resin by blending an acrylic resin, such as polymethylmethacrylate (PMMA), with the polyacetal as a base resin, merely blending the acrylic resin with polyacetal has its own problems. That is, a blend of polyacetal resin and acrylic resin typically cannot be utilized as an exterior component (especially components that are large or complicated in shape) because the characteristic mechanical properties of the polyacetal resin (such as its impact resistance and tensile elongation properties) are impaired by the presence of the polymethylmethacrylate. Although it is not known precisely why such polyacetal property deterioration occurs, it is surmised that it probably occurs due to the poor dispersion of the polymethylmethacrylate in the polyacetal base resin. Moreover, blending polymethylmethacrylate resin with polyacetal base resins does not improve the dimensional accuracy, such as distortion and shrinkage during molding, to sufficient degrees. It is therefore towards solving these problems that the present invention is directed.\nThe present invention is broadly directed to polyacetal resin compositions and processes for making the same whereby a polyacetal base resin, a particulate filler material and an acrylic resin may be blended with one another under melt conditions and in certain specified amounts so as to improve the dispersion of the acrylic resin and filler throughout the polyacetal base resin. The present invention more specifically relates to polyacetal compositions whereby an interconnected network consisting essentially of a continuous phase of the acrylic resin is formed within the polyacetal base resin (which therefore forms a matrix for the continuous acrylic resin phase). The particulate filler material is dispersed homogeneously throughout the polyacetal base resin but is encapsulated within the continuous acrylic resin phase. In this connection, it is most preferred that the filler material exhibits a greater surface tension at melt-blending temperatures as compared to the surface tension of the acrylic resin forming the continuous phase in the polyacetal matrix base resin.\nFurther aspects and advantages of this invention will become more clear from the following detailed description of the preferred exemplary embodiments thereof."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of document searching, data mining and data visualization.\n2. Description of the Related Art\nThe field of data searching and data/text mining is replete with various search methods and algorithms for helping determine the identity and/or location of documents that may have relevance to a particular subject matter of interest. The most basic search techniques involve locating specific words or word combinations within one or more of a quantity of documents contained in a database. This search methodology, while very simple to implement, suffers from a number of significant drawbacks, including slow search processing time, limited ability to construct and execute complex search queries, and other well-documented limitations inherent in the use of keywords as search criteria. Improvements to the basic keyword search include the use of structured queries (e.g., based on Boolean logic), word stemming, wildcards, fuzzy logic, contextual analysis and latent semantic analysis.\nDespite its well-documented drawbacks, simple key-word based searching is still a good entry point to quickly locate documents of general interest to a relevant subject matter. It is sufficient in many searching applications to locate a particular desired piece of information contained within one or more documents being searched. However, there are many specialized searching applications, particularly in the science, technology, academic and legal fields, where keyword searching (even with the various improvements to date) provides an unsatisfactory approach for locating some or all of the relevant documents that may be of interest to a researcher. The primary underlying difficulty is that words and word phrases are imprecise by their nature. Different words and word phrases can have completely different meanings in different associative contexts. As a result, key-word based searching in these and other specialized searching applications tends to be a slow and tedious process, typically producing significant numbers of irrelevant documents or “false hits” and often failing to turn up one or more desired relevant documents.\nMore advanced searching techniques rely on contextual or bibliographical linkages between two or more documents. For example, U.S. Pat. No. 6,754,873 issued Jun. 22, 2004 to Law, et. al. describes a search technique for finding related hyperlinked documents located on the world-wide-web using link-based analysis. In this case backlink and forwardlink sets are utilized to find web pages that are related to a particular selected web page of interest. The resulting list of related web pages is typically sorted in accordance with a calculated relevancy score, the intent being that presumably the most relevant and/or highest quality hits would be listed toward the top of the search results page and the least relevant and/or lowest quality hits would be listed toward the bottom of the search results page.\nRelevancy scores are typically calculated as an arbitrary score or metric based on one or more selected factors determined (or assumed) to be informative as to the quality or relevance of the search output relative to the search input. For example, the search engine may assign an arbitrary rank or score to each hit calculated according to the number or frequency of keyword occurrences in each document, the intent being that the total score would roughly correspond to the relevance or importance of the particular located document relative to the input search query. Another example, described in the article entitled “The Anatomy of a Large-Scale Hypertextual Search Engine,” by Sergey Brin and Lawrence Page, assigns a degree of importance to a web page based on the link structure of the web page. In this manner, the Brin and Page algorithm attempts to quantify the importance of a web page based not on its content, but on the number and quality of linkages to and from other web pages.\nU.S. Pat. No. 6,526,440 issued Feb. 25, 2003 to Bharat and assigned to Google, Inc. describes a similar search engine for searching a corpus of data and refining a standard relevancy score based on the interconnectivity of the initially returned set of documents. The search engine obtains an initial set of relevant documents by matching search terms to an index of a corpus. A re-ranking component in the search engine then refines the initially returned document rankings so that documents frequently cited in the initial set of relevant documents are preferred over documents that are less frequently cited within the initial set. The resulting hits in each case are typically displayed in a text-scrolled list, with the relative placement of each hit on the list being determined in accordance with the calculated relevancy score. This, in essence, is the primary search and relevance ranking algorithm behind the popular Google® search engine.\nAs with the Google® search engine, many of the more sophisticated search engines today are primarily optimized toward the task of searching the world wide web for relevant documents of a general-content nature and focusing typically on a single item of information or a single concept. Most searches conducted using these types of search algorithms seek to find particular items of information that are essentially known to exist and that can be described with a few simple key words. The probability that a user would be able to successfully use a search engine in this context to locate at least one source of information satisfying the user's need is fairly high. However, in certain specialized searching applications, particularly in the science, technology, academic and legal fields, conventional search engines provide an unsatisfactory approach for locating some or all of the relevant documents that may be of interest to a researcher.\nFor example, those skilled in the intellectual property arts and the patent legal field in general will readily appreciate the difficulty and challenge of searching through vast databases of case law, patents and related scientific documents looking for “prior art” documents relevant to a particular issued patent or pending application and/or cases relevant to a particular point of law. For patents the difficulty and challenge stems from the confluence of several unique factors affecting patents and patent-related documents. These factors include the shear volume of potentially relevant patent documents and related scientific literature (estimated at over 80 million documents worldwide), latent inaccuracies and inconsistencies in the technology classifications used by the various national and international patent offices, the complex scientific nature of patent disclosures, the ever evolving lexicon for describing novel patented concepts and structures, language translation issues in the case of relevant foreign patent documents and scientific literature, and the proclivity of patent attorneys and agents to use complex legalese and coined lexicon to describe novel concepts. The purpose of the patent search is also quite different than the normal search context. The point is not so much to find useful information relevant to a concept of interest, but to establish and document legal evidence of the existence or non-existence of a particular concept or idea in combination with one or more other related concepts or ideas at a particular point in time.\nTraditional search engines are not particularly adept at efficiently handling these and other types of specialized searching applications. The standard input/output text interface of most conventional search engines also does a poor job of displaying and communicating input/output search criteria and search results in a way that facilitates intuitive understanding and visualization of the logical relationships sought to be explored between two or more related concepts being searched. It would be of particular benefit to provide an improved search algorithm, database and user interface that would overcome or at least mitigate some or all of the above-noted problems and limitations."}
{"text": "A wide variety of equipment and systems, such as portable and desktop computers, mainframe computers, communication infrastructure frames, automotive equipment, etc., include heat-generating components in their casings. One method used to remove heat from such equipment is to have a tube-axial fan draw air from exterior of the housing to blow cooling air over the heat-generating components. It has also been the trend to make electronic equipment more compact with the casings or housings therefor becoming smaller for consumer use and transport with highly complex circuits and miniaturized integrated electronic components used therewith, creating high-density electronic systems in the small casings. Thus, the heat generated in the casing can become relatively intense and requires high performance from the fan to meet the additional cooling demand.\nA known problem that axial fans face relates to the amount of airflow that can be generated by the fan before the system resistance creates a static pressure drop or stall effect, especially in high-density applications. This typically is shown in a static pressure (P) vs. quantity of air flow (Q) graph plotting fan performance and system impedance pressure curves thereon. Greater heat load systems, such as the aforedescribed high density electronic systems, require a larger amount of airflow against high pressure drop to provide the desired cooling for the system. In such applications, this increased cooling demand is typically not available in an efficient manner from a conventional fan because of the relatively early stalling effect of the fan when it cannot overcome the pressure increasing in the system as it is being cooled.\nMore specifically, it is the section of the fan performance curve between the point where the static pressure is zero, and the point where the static pressure is at its peak value and before system pressure begins to reduce the static pressure from its maximum, which is the efficient or preferred operating area of the fan and which is desired to be maximized in terms of fan performance. At the high pressure end of the fan preferred operating area, the fan operation is in a \"centrex region\" which is an area proximate and at the intersection of both the fan performance and the system pressure curves. In the centrex region, as the output pressure from the fan drops with decreasing flow rate, the fan begins to show stall characteristics. Fan operation in the centrex region is in a highly unstable and inefficient manner. As flow reduces and output pressure increases after the centrex region, the fan operates in a highly turbulent region and mixed mode of operation which can increase vibrations and power requirements for the fan.\nA known method to meet higher heat loads is to employ multiple fans aligned in series in a push/pull configuration. The series multiple-fan system has the ability to raise the maximum pressure before stall. The shortcomings with series multiple fans include the increased space necessary for incorporating multiple fans into the smaller casings of high density electronic system and the additional power requirements. Space requirements in electronic devices such as microsite communication infrastructure systems are of paramount importance, and accordingly space allocation for an additional fan may not be feasible. Multiple fans also produce an increased level of noise which is undesirable, particularly where the fans are incorporated into consumer devices. At the bottom line, the cost associated with multiple fans is also increased.\nU.S. Pat. No. 5,707,205 discloses a fan with a special construction including annular plates disposed about the fan blades and spaced axially by an optimum amount to improve flow. The fan of the '205 patent with the annular plate construction is disclosed as improving the so-called P-Q characteristic of the fan but has been found to lack the desired airflow in the preferred operating area of the fan (i.e., from zero static pressure to the maximum pressure before stall) and specifically at lower pressures thereof because of the oversized nature of the circumferentially extending spaces between adjacent plates.\nThus, there is a need for an improved air moving device, such as a fan used to cool high-density electronic systems, that has an increased maximum pressure before the stall effect. Further, a fan would be desirable that does not reduce airflow in the preferred operating area. Also, it is desired that the fan be able to be incorporated into the smaller housing including high density electronic systems therein in a cost-effective manner and without creating excessive noise."}
{"text": "The present invention relates to a glass-ceramic panel providing a cooking surface with at least one cooking zone, which is associated with a heating element with an annular or ring-shaped lighting device and which is marked off by an optical structured cooking zone indicating device visible on the cooking surface.\nThe present invention also relates to a method of making the glass-ceramic panel providing the cooking surface with the at least one cooking zone, which is associated with the heating element and provided with the structured cooking zone indicating device.\nGlass-ceramic panels providing cooking surfaces with individual cooking zones are established parts of modern kitchen equipment. The cooking zones typically are part of the glass-ceramic panel. However they can also be formed by separate ceramic or metallic parts, which are inserted in openings in the glass-ceramic panel. The cooking zones on the cooking surface are marked off optically in order to make the manipulations of placing cooking or roasting vessels on the respective cooking zones easier and for optimum usage of heat energy.\nA number of methods have been used in order to produce this sort of cooking zone marking off device, which is designated in the following as a cooking zone indicating device.\nFor example, it is possible to make it by providing a suitable pattern on the cooking surface that has a different appearance than the remaining surface of the glass-ceramic panel by means of a suitable treatment during the course of the manufacture of the glass-ceramic panel.\nIn order to guarantee a uniform appearance for the cooking surface portions that do not operate as cooking areas, numerous optical arrangements are known for producing the cooking zone indicating device under the glass-ceramic panel. They typically have lighting elements connected with light guides in different configurations, in order both to produce a through-going and also point-illuminating cooking zone indicating device, if necessary, in combination with a residual heat indicator.\nThese optical display elements are associated directly with the heating element arranged under the cooking zone, which typically is formed by a radiating heated body, but can also be embodied by an inductively acting heating element. For example, EP 0 540 876 A2 describes an externally lighted glass ring formed as a light conductor, which is integrated in an insulating ring of the heating element. This illuminated glass ring indicates the circular cooking zone inside of the ring and can be used also as a residual heat indicator.\nAn annular region marked off on the cooking surface acting as cooking zone indicating device is described in DE 40 04 309 A1. This cooking zone indicating device marks off this annular region in contrast to other parts of the cooking surface even more strongly by design elements.\nDE 197 00 551 A1 discloses an annular state indicating device in the form of a lighted ring, which is arranged around a conventional heated cooking plate embedded in a glass-ceramic panel. An annular transparent region of the glass-ceramic panel is formed by a suitable annular interruption in a non-transparent cover layer on the underside of the glass-ceramic panel.\nNo structure is provided in the annular indicating regions in both these known cooking zone indicating devices.\nA structured lighted cooking zone indicating device is described in EP 0 746 179 A2, which is equivalent to U.S. Pat. No. 5,750,959. In this cooking zone indicating device the previously described illuminated glass ring is covered by a light-blocking or light-impermeable heat-resistant foil provided with holes and arranged between the insulating wall for the heating element and the glass-ceramic underside, in order to provide a structured cooking zone indicating device. Since this foil is preferably mica and is only structured by a comparatively difficult cutting method, e.g. laser cutting methods, this method of providing a structured lighted cooking zone indicating device is comparatively expensive.\nFurther correct placement of this cooking zone indicating device during assembly in the cooking area is very troublesome, time-consuming and expensive.\nStructured cooking zone indicating devices are currently in fashion in comparison to through-going cooking zone indicating devices and are frequently demanded by customers.\nIt is an object of the present invention to provide a glass-ceramic panel of the above-described type that provides a cooking surface with at least one cooking zone, which is associated with a heating element with an annular lighting device and which is marked off optically by a structured cooking zone indicating device on the cooking surface by means of the annular lighting device, by simple means, which is more economical and faster to make than the corresponding known mica mask and which guarantees easier handling during assembly.\nIt is also an object of the present invention to provide a method of making the above-described glass-ceramic panel providing a cooking surface with at least one cooking zone, which is associated with a heating element having an annular lighting device and which is marked off optically by a structured cooking zone indicating device by means of the lighting device, wherein this method is more economical and faster than the current method and which guarantees easier handling during assembly.\nThese objects, and others that will be made more apparent hereinafter, are attained in a glass-ceramic panel providing a cooking surface with at least one cooking zone, which is associated with a heating element having an annular lighting device and which is marked off optically by a structured cooking zone indicating device by means of the annular lighting device.\nAccording to the invention a structured light-blocking or light-impenetrable covering layer is applied directly to the underside of the glass-ceramic panel and provided with transparent regions within the structured cooking zone indicating device corresponding to or according to a desired predetermined pattern or structure for the cooking zone indicating device.\nThe application of the coating layer with the structured light-permeable regions on the underside of the glass-ceramic panel can occur economically and faster than the production of a mica mask. Since the structure fixed to the underside of the glass-ceramic panel corresponds to the later position of the heating element in the cooking area, the assembly of the glass-ceramic panel requires no special handling in regard to the cooking zone indicating device. Since these known predetermined light-permeable regions are provided by application of the structure pattern, the light from the annular lighting device, preferably in the form of a lighted glass ring, passes through these holes or openings in the coating layer, thus forming the cooking zone indicating device.\nThe method of making the glass-ceramic panel providing the cooking surface with the at least one cooking zone, which is associated with the heating body and which has the structured cooking zone indicating device, includes the steps of:\nprinting the coating layer made from heat-resistant light-impermeable material provided with transparent regions corresponding to the desired or predetermined structure or pattern within the cooking zone indicating device directly on the underside of the glass-ceramic panel by means of screen printing; and\ndrying the resulting printed glass-ceramic panel.\nThis screen printing technique guarantees an economical and fast application of the structured coating layer on the underside of the glass-ceramic panel with a minimum of time consumed.\nGlass-ceramic cooking surfaces are known, which are light-impermeable in an unheated region from the underside. The inner regions of the cooking area are covered by illumination from the outside. For that purpose a smooth or flat organic light-impermeable enamel coating is applied to the underside of the glass-ceramic cooking surface in the appropriate regions by means of screen printing.\nIn the known embodiment the flat underside coating only serves for covering the unheated regions. There is no disclosure in the prior art that suggests the structured coating for designating the heated cooking zone and for conveying information regarding it to a user."}
{"text": "The present invention relates to database systems.\nIn conventional relational databases, all data are stored in named tables. The tables are described by their features. In other words, the rows of each table contain items of identical type, and the definitions of the columns of the table (i.e., the column names and the data types stored in the column) describe the attributes of each of the instances of the object. By identifying its name, its column names and the data types of the column contents, a table is completely described. Queries to a relational database are formulated in a query language. One such language is SQL (Structure Query Language) which is widely used in commercial relational database systems. The data types offered by SQL can be classified as character arrays (names), numbers, and data types related to date and time. Tables can be modified or combined by several operations of relational algebra such as the application of Boolean operators, projection (i.e. selection of columns) or the Cartesian product.\nRelational databases offer several advantages. Database queries are based on a comparison of the table contents. Thus, no pointers are required in relational databases, and all relations are treated uniformly. Further, the tables are independent (they are not related by pointers), so it is easier to maintain dynamic data sets. The tables are easily expandable by simply adding new columns. Also, it is relatively easy to create user-specific views from relational databases.\nThere are, however, a number of disadvantages associated with relational databases as well. For example, access to data by reference to properties is not optimal in the classical relational data model. This can make such databases cumbersome in many applications.\nAnother recent technology for database systems is referred to as object oriented database systems. These systems offer more complex data types in order to overcome the restrictions of conventional relational databases. In the context of object oriented database models, an “object” includes both data and the functions (or methods) which can be applied to the object. Each object is a concrete instance of an object class defining the attributes and methods of all its instances. Each instance has its unique identifier by which it can be referred to in the database.\nObject oriented databases operate under a number of principles. One such principle is referred to as inheritance. Inheritance means that new object classes can be derived from another class. The new classes inherit the attributes and methods of the other class (the super-class) and offer additional attributes and operations. An instance of the derived class is also an instance of the super-class. Therefore, the relation between a derived class and its super-class is referred to as the “is A” relation.\nA second principle related to object oriented databases is referred to as “aggregation.” Aggregation means that composite objects may be constructed as consisting of a set of elementary objects. A “container object” can communicate with the objects contained therein by their methods of the contained objects. The relation between the container object and its components is called a “partof” relation because a component is a part of the container object.\nYet another principle related to object oriented databases is referred to as encapsulation. According to encapsulation, an application can only communicate with an object through messages. The operations provided by an object define the set of messages which can be understood by the object. No other operations can be applied to the object.\nAnother principle related to object oriented databases is referred to as polymorphism. Polymorphism means that derived classes may re-define methods of their super-classes.\nObjects present a variety of advantages. For example, operations are an important part of objects. Because the implementations of the operations are hidden to an application, objects can be more easily used by application programs. Further, an object class can be provided as an abstract description for a wide variety of actual objects, and new classes can be derived from the base class. Thus, if an application knows the abstract description and using only the methods provided by, the application can still accommodate objects of the derived classes, because the objects in the derived classes inherit these methods. However, object oriented databases are not yet as widely used in commercial products as relational databases.\nYet another database technology attempts to combine the advantages of the wide acceptance of relational databases and the benefits of the object oriented paradigm. This technology is referred to as object-relational database systems. These databases employ a data model that attempts to add object oriented characteristics to tables. All persistent (database) information is still in tables, but some of the tabular entries can have richer data structure. These data structures are referred to as abstract data types (ADTs). An ADT is a data type that is constructed by combining basic alphanumeric data types. The support for abstract data types presents certain advantages. For example, the methods associated with the new data type can be used to index, store, and retrieve records based on the content of the new data type.\nSome conventional object-relational databases support an extended form of SQL, sometimes referred to as ObjectSQL. The extensions are provided to support the object model (e.g., queries involving object attributes). However, these object-relational databases are still relational because the data is stored in tables of rows and columns, and SQL, with some extensions, is the language for data definition, manipulation, and query. Both the target of a query and the result of a query are still tables. The extended SQL language is often still the primary interface to the database. Therefore, there is no direct support of host object languages and their objects. This forces programmers to continue to translate between objects and tables.\nData pertaining to the operation of a business such as types of companies forming an enterprise, orders that the company receives from various customers, and what constitutes an order is hierarchical. As discussed above, relational systems for storing data, on the other hand, are tabular in nature, and consequently, do not directly represent hierarchies. Accordingly, object programming models do not represent the hierarchy of business data very well."}
{"text": "US 2008/0080770 A1 discloses a method and a system for visualizing regions in an image data set, and, in particular, discloses a method for assigning labels to regions in an image data set. The method comprises the steps of deriving a probabilistic model for a plurality of geometrical structures, computing a regional response around a region in the image data set, computing a region score for each geometrical structure using the plurality of probabilistic models and labeling the region in the image data set based on the region score. In a particular embodiment, the principle curvatures of iso-surfaces at each voxel of the image data set are used as the regional response and the region score represents the probability of observing the computed curvature data for a region corresponding to a specific anatomical structure, such as a vessel or a nodule. The labeled regions may then be displayed to a radiologist with different colors for different anatomical structures.\nThis visualization of the image data set needs at lot of computational intensive calculations like the derivation of a plurality of probabilistic models for a plurality of geometrical structures, computing a regional response around a region in the image data set and computing a region score for each geometrical structure using the plurality of probabilistic models. The visualization of the image data set such that different objects are visually separated from each other requires therefore large computational costs."}
{"text": "A variety of electric squibs have been developed as squibs for gas generators for inflating air bags equipped in cars and squibs for micro gas-generators used for seat belt pretensioners.\nSuch a squib usually has metal pins for electrically connecting to the external, and a heating element at the other ends of the metal pins for igniting an explosive.\nAs examples of such a heating element, heating elements incorporated in a printed sub-circuit have been known as disclosed in a patent document 1 and a patent document 2 corresponding to the patent document 1.\nOn the other hand, moreover, it has been contemplated to link an air bag system to a local area network (LAN) to control ignition of a squib by communication.\nIn this case, it is necessary to locate electric circuits in the squib for the communication and ignition as described in a patent document 3.\nWith the squib constructed as described above, therefore, it becomes necessary to provide means for causing electric current to flow through the electric circuits built in the squib.\nIn a patent document 4, for example, disclosed is a structure of an electric circuit board fixed to electrode pins of a closure plug by soldering or the like.\nIn the method disclosed in the patent document 4, after the circuits have been previously constructed on a board and connected to electrode pins of a closure plug, a resin must be molded around the circuits and the board.\nHowever, particularly the resin molding produced under a condition already connected to a closure plug is seriously lower in productivity as compared with the molding of parts such as a usual IC fixed to a lead frame.\nMoreover, in the case that connection between the electrode pins and the electric circuits is firmly fixed by soldering, there would be a risk of solder cracks due to residual stresses.    Patent document 1: Specification of French Patent Application Laid Open No. 2,704,944    Patent document 2: Specification of U.S. Pat. No. 5,544,585    Patent document 3: Official Gazette of Japanese Patent No. 3,294,582    Patent document 4: Specification of European Patent Application Laid Open No. 1,256,775"}
{"text": "Magnetic Resonance Imaging (MRI) is a medical imaging technique commonly used to image inside the human body. Generally a strong magnetic field (e.g., 1.5 T) is applied to the patient to align the nuclear magnetization vectors of hydrogen atoms in the water of the tissue of the patient under investigation. Simultaneously a RF (e.g., 63.86 MHz) magnetic field (˜0.14 μT) is applied to perturb the aligned magnetization, causing the hydrogen nuclei to emit energy signatures that are detectable by a scanner in the MRI system.\nThe MRI signal is used to construct an image. Different tissues are detected because the protons in various tissues return to their equilibrium nuclear magnetization states at different rates. This effect is used to create contrast among different types of tissues. Diseased tissues, such as tumors, can be detected in this manner as well.\nA time-varying magnetic field creates an electric field within any electrically conductive material and the electric field, in turn, induces an electric current, referred to as the “eddy current”. Generally the electromagnetic field penetrates into metals up to a certain depth, commonly referenced to a “skin depth”. An electromagnetic wave entering a conducting surface is damped so that the current density is largest near the surface of the conductor and reduces in amplitude by a factor 1/e at a distance referred to as the skin depth (δ) of the electrical conductor from its surface given by:δ=[2/(ωμ0σ)]1/2 where ω is the angular frequency of the radiation, μ0 is the permeability of the radiation in a vacuum, and σ is the electrical conductivity of the metal. δ can be seen to decrease if the electrical conductivity (σ) increases. For example, at ω=60 Hz in copper, δ is about 8.5 mm. At high frequencies δ may be much smaller (shallower).\nThe eddy current can heat up the electrically conductive material by the Joule effect and this process is called induction heating. Induction heating can cause thermal damage to tissues while conducting MRI scans for patients wearing implants that have leads, such as in spinal fusion stimulators, cardiac pacemakers and neurostimulation systems. A maximum temperature change of 25.38° C. has been reported for a deep brain stimulation implant shortly after initiating MRI. The heating can be more severe at metal tips. A temperature elevation of 63.18° C. has been reported to have occurred at the tip of a pacing electrode (unattached to a cardiac pacing pulse generator) within 90 s of initiating MRI.\nAlthough selection of a conventional high electrical conductivity metal material can reduce MRI-induced induction heating somewhat by reducing the amount of energy that is actually introduced into the material via the skin effect, the resulting MRI-induced induction heating may still be too high for certain applications. There is thus a need for new materials and enabling processing systems and methods for forming such materials which provide modified surfaces which generate reduced MRI-induced induction heating."}
{"text": "The present invention relates to an aparatus for restoring the cell storage regions of an ATM (Asynchronous Transfer Mode) switch having a common buffer structure, and a method therefor.\nGenerally, data is communicated as cells in an ATM network, and each of the cells consists of 53 bytes including a header of 5 bytes and a payload of 48 bytes. The header stores information of an address, a sender, etc., while the payload stores content data in order to distribute such data cells according to addresses, there may be employed an ATM switch having a common memory structure. The earlier ATM switch comprises a first multiplexer for multiplexing the cells with the header and content data separated, a common memory for storing the content data generated from the first multiplexer, a plurality of FIFO (first-in first-out) registers having the same number as the output ports to store the addresses for the content data by the headers, an idle addresses for the content data stored in the common memory, a second multiplexer for multiplexing the outputs of the FIFO registers, an address checker for checking the state of the addresses generated from the second multiplexer to store them into the idle address pool, and a demultiplexer for distributing the data output from the common memory among the output ports.\nIn operation, an ATM cell applied from a subscriber\"\"s board or another system to the ATM switch is separated by the multiplexer into the header and payload. The header determines the output port for the content data stored in the payload of the ATM cell. The content data is stored into the common memory at an address designated by the idle address pool storing the idle addresses of the common memory. The header is stored into the corresponding FIFO register together with the address where the content data is stored into the common memory. Each register generates the address where the header and content data of the ATM cell are stored towards the multiplexer according to a read signal generated from a read time generator. Then the multiplexer generates towards the address checker the address of the common memory where the content data of the ATM cell, and towards the common memory a read signal to read the data of the address transferred to the address checker. The data read from the common memory is demultiplexed by the demultiplexer generated through respective output port. Meanwhile, the address checker checks if there is an error in the received addresses. The address checker cancels the erroneous address, or transfers the correct address to the idle address pool (IAP). The IAP goes on storing the addresses where the content data may be stored into the common memory.\nIn this case, whenever the address checker cancels erroneous address, there is reduced the number of the addresses where the content data may be stored into the common memory. Namely, although there has occurred to errors in the common memory, the addresses where the content data may be stored into the common memory are converted into the error state, thus making it impossible to use the storage regions of the erroneous addresses. If stored in the IAP, the erroneous address stored into the IAP may cause a loss of the data of the ATM cell. However, if allowed to continue, the number of erroneous addresses stored into the address checker from the registers through the multiplexer, the number of the addressses stored in the IAP becomes reduced accordingly. In order to restore the storage regions of the common memory, the operator must detach the board mounted with the switch from or reset the system to restore the IAP.\nTakechi et al U.S. Pat. No. 5,513,191 discloses an Asynchronous Transfer Mode (ATM) Cell Error Processing System. In addition to an error detector, Takechi et al discloses an error editing unit. This error editing unit has the function of determining for each of the ATM cells, whether based on the respective and corresponding decision signals, a respective ATM cell should be discarded. Takechi et al does not indicate that the error editing unit restores error cells because of an erroneous address. What is needed is a method and apparatus for ATM networks for restoring cells having erroneous addresses.\nIt is an object of the present invention to provide an apparatus for restoring the idle addresses of the common memory in an ATM switch system without detaching or resetting the system.\nIt is yet another object of the present invention to provide a method and apparatus for an ATM network for restoring cells that have erroneous addresses.\nAccording to one aspect of the present invention, an apparatus for restoring the cell data storage regions of the common memory in an asynchronous transfer mode (ATM) switch system, comprises a first multiplexer for multiplexing the cell data with the header and content data separated, a second multiplexer for multiplexing the addresses generated from a plurality of FIFO registers, an address checker for checking the addresses from the second multiplexer to generate a first and second address signals if the addresses are checked normal, an address memory storing all addresses of the common memory in error state to selectively convert the addresses into normal state according to the second address signal from the address checker, a controller for checking the address memory at predetermined intervals to restore an address stored in error state to a normal state, and an address multiplexer for generating an idle address to an IAP according to the first address signal and address restored signal.\nAccording to another aspect of the present invention, an ATM switch system, which includes a common memory, an address checker, address multiplexer and controller for checking the address memory to restore the addresses of the common memory, is provided with a method for restoring the cell data storage regions of the common memory, which comprises the steps of causing the address checker to check during operation of the timer of the controller whether multiplexed read addresses are correct, writing corresponding storage regions of the address memory with correct addresses or canceling erroneous addresses according as the multiplexed read addresses are correct or not, sequentially checking the addresses of the address memory upon termination of the timer, and restoring the erroneous addresses checked in the previous step into correct state,"}
{"text": "The touch control technology used in touch panels can be classified into resistive type, capacitive type, surface acoustic wave type, and optics type, based on their sensing principles.\nFIGS. 1 and 2 are a block diagram and a circuit diagram of a control circuit of a conventional capacitive touch panel, respectively. The control circuit 10 is electrically connected to the capacitive touch panel 12, and comprises a pulse generator 20, a multiplexer 22, a multiplexer 24, an integrator 26, a sampling and holding circuit 28, and an analog to digital converter (hereinafter called the ADC) 32.\nThe pulse generated by the pulse generator 20 functions as a driving which is transmitted through the multiplexer 22 to the X-axial driving wires and the Y-axial driving wires of the capacitive touch panel 12, and the value of capacitive change sensed by the Y-axial and X-axial sensing wires of the capacitive touch panel 12 functions as a touch control sensing signal, which is transmitted through the multiplexer 24 to the integrator 26.\nIn FIG. 2, the touch control sensing signal of the capacitive touch panel 12 transmitted through the multiplexer 24 to the integrator 26 is a pulse signal, and the integrator 26 integrates several pulse signals into a ladder-shaped wave signal and transmits it to a sampling and holding circuit 28.\nThen, in FIG. 1, the sampling and holding circuit 28 samples the ladder-shaped wave signal. integrated by the integrator 26, and keeps the sampled signal at a direct current (DC) voltage level. The ADC 32 will convert the aforementioned sampled analog signal into a digital signal and transmit the digital signal to a microprocessor (not shown) to conduct a signal processing.\nThe aforementioned control circuit 10 used in the conventional capacitive touch panel 12, at least, has the following disadvantages:\nThere is noise signal generated in the circuit environment, and the noise signal will also be integrated by the integrator 26 as a portion of the signal, resulting in a low signal to noise ratio (SNR) of the whole system.\nIt takes a period of time for the integrator 26 to transfer and accumulate the charges on the capacitive touch panel 12 to a voltage that can be processed by a rear-end circuit. Such a long integration time will affect the frame rate of the capacitive touch panel 12. This problem may be compensated by way of a specific hardware (for example, the sampling and holding circuit 28); however, the cost will increase.\nThe use of the integrator 26, sampling and holding circuit 28, and digital to analog converter 32, will occupy a larger area on the chip."}
{"text": "One type of implantable device is a synthetic vascular graft such as is commonly used to replace damaged or dysfunctional arterial or venous pathways, for example at the site of an aneurysm or occlusion. Bypass grafts are often used to divert blood flow around damaged regions to restore blood flow. Another use of vascular prostheses is for creating a bypass shunt between an artery and vein, specifically for multiple needle access, such as is required for hemodialysis treatments. Following multiple percutaneous invasions into a vein, the vein may either collapse along the puncture track or become aneurysmal, leaky or fill with clot, causing significant risk of pulmonary embolization. Vascular prostheses have been used for many years as an alternative to patients' own veins for vascular access during hemodialysis.\nMaterials research has led to the development of some synthetic materials for use in artificial vascular prostheses. For example, polytetrafluoroethylene (PTFE), a polymeric material which may be stretched to a specific length and expanded to a specific thickness, is often used to fabricate single lumen artificial veins and arteries. Typically, however, PTFE vascular grafts cannot safely be used to withdraw blood until they have been in place in the body for a minimum of 14 days after surgery and have become surrounded by fibrotic tissue. This is because bleeding occurs at the site of a needle puncture in PTFE grafts if fibrotic tissue is absent. Complications which can result from early puncturing of PTFE arteriovenous fistulas include a hematoma surrounding the graft, false aneurysm, and graft occlusion.\nVarious other synthetic materials, in addition to PTFE, have been used for vascular grafts, including Dacron.RTM. brand and other synthetic polyester fibers, mandrel spun polyurethane, and silicon elastomer fibers. Additionally, vascular grafts have been formed using autologous saphenous vein, modified bovine carotid xenograft, and modified human umbilical vein. None, however, has overcome the problems associated with early failure of the graft following implantation.\nA problem associated with known vascular grafts stems from their delicate structure making them difficult to handle and position properly during surgery. Due to their circumferentially uniform appearance, the grafts may be twisted during implantation, which can reduce the openness, or patency, of the implanted graft.\nAnother problem associated with known vascular grafts is that they are transparent to known non-surgical techniques for viewing that are generally used to detect structures in the body. These techniques include x-rays, MRI scanning, fluoroscopy, ultrasound, and nuclear magnetic resonance. As a result, post-implantation examination of known vascular grafts is difficult. If a surgeon suspects that a blockage of an implanted vascular graft has occurred, for example, she must inject radiopaque dye into the patient through the graft. The dye allows the graft to be visible during fluoroscopic examination to determine whether it has collapsed or is satisfactorily transporting blood flow. This process is invasive to the patient and places further burdens on the patient's circulatory system.\nVarious marking devices have been developed to avoid the need for injecting dye into a patient who has undergone vascular graft implantation. U.S. Pat. No. 4,202,349 (Jones, May 13, 1980), for example, discloses a radiopaque blood vessel marker for attachment to the side wall of a blood vessel or vascular graft. The marker is a substantially flattened disk which is sutured into place by the surgeon. The Jones vessel markers are not well suited for use with PTFE vascular prostheses because it is difficult to elicit natural occlusion of suture holes in PTFE due to the relative elasticity of the porous PTFE tubing materials.\nU.S. Pat. No. 4,787,391 (Elefteriades Nov. 29, 1988) discloses a device for the marking the anastomosis, or joining site, of a coronary graft so that the anastomosis can be localized fluoroscopically after surgery. The device includes gold indicators which are arranged on a hemostatic material which is absorbed or dissolved completely in the body during the healing period, leaving the gold markers at the site. This device is useful for marking the area of the anastomosis, but does not allow the vascular graft itself to be viewed during fluoroscopic examination.\nIt is, therefore, an object of the present invention to provide a vascular graft having integral diagnostic indicia for determining the patency (openness) of the graft vessel by non-invasive techniques. It is also an object of the present invention to provide a vascular graft having integral radiopaque or MRI markers that do not need to be sutured to the graft. It is still another object of the present invention to provide a method for treating a patient having a damaged or dysfunctional arterial or venous pathway using the detectible vascular grafts of the present invention."}
{"text": "MYLOTARG® (gemtuzumab ozogamicin), also referred to as CMA-676 or just CMA, consists of a monoclonal antibody against CD33 that is bound to calicheamicin by means of an acid-hydrolyzable linker. When the derivatized calicheamicin binds to the DNA minor groove, it disrupts DNA progression and eventually causes cancer cell death. The commercial product was marketed as the first antibody-targeted chemotherapeutic agent under the name MYLOTARG® and is currently approved for the treatment of acute myeloid leukemia (AML) in elderly patients.\n3-Methyl-3-mercaptobutanoic acid hydrazide, also called DMH linker, or CL-332258, is used to link calicheamicin to monoclonal antibodies. The derivatized calicheamicin is then activated for conjugation with a humanized monoclonal antibody to give CMA-676. Currently, DMH linker is prepared through a 5-step reaction process via the intermediate, p-methoxybenzylthioether hydrazide, 5. (Equations I-V). In the present manufacturing process, a Michael addition of p-methoxy-benzylthiol to 3,3 dimethylacrylic acid is assisted by piperidine, (Equation I).\n\nThe resulting thioether acid (1) reacts with oxalyl chloride in methylene chloride to form p-methoxybenzylthioether acid chloride (2) (Equation II).\nAcid chloride (2) is slowly added to a mixture of hydrazine/methylene chloride (in a ratio of about 28%, v/v) at low temperature (−70° C.). The corresponding p-methoxybenzylthioether hydrazide (3) forms in about 74% yield (Equation III):\n\nHowever, the desired product p-methoxybenzylthioether hydrazide (3) typically contains about 20% of an undesired by-product, bis-methoxybenzylthioether hydrazide (see Equation VI below). Removing the benzyl protecting group under acidic conditions (Equation (IV), followed by neutralization of the acid salt and purification (Equation V) provides 3-methyl-3-mercaptobutanoic acid hydrazide (5) in 45% yield.\n\nAn undesired by-product, bis-methoxybenzylthioether hydrazide (6) is generated from the reaction of the product p-methoxybenzylthioether hydrazide with the starting material p-methoxybenzylthioether acid chloride (Equation VI). The generation of this by-product results in lower yield and quality.\n\nUsing original process procedures, the bis-methoxybenzylthioether hydrazide (6), is generated in amounts of about 20%. The presence of this level or greater of undesired by-product from Equation III is clearly undesirable. The present invention describes techniques which overcome this problem and decrease the yield of the undesired by-product."}
{"text": "The invention relates to a sketching aid which facilitates the sketching or drawing of a real-life scene, allowing duplication of proper perspective and proper spatial relationships between component parts of the scene to be sketched. Utilizing the present system, a skilled artist or draftsman can much more quickly, and with good accuracy, construct a perspective view of a scene or object. For a relatively unskilled artist or draftsman, it becomes possible, utilizing the sketching aid according to the present invention, to make a sketch or drawing having a much higher degree of realism than would otherwise be possible. The sketching aid according to the invention can be used by a wide variety of people, both for enjoyment or in the normal course of their profession, or in accordance with a course of study.\nThe sketching aid according to the present invention is constructed so that it is easy to handle and transport, is versatile, is complete, and allows the basic functions for which it is designed to be practiced readily by any user. The device according to the invention provides for storage of transparent drawing films directly with it, appropriate mounting of a drawing implement with the sketching aid, and comfortable and functional support of the aid by a single hand of the user, at any desired relative position with respect to the user.\nAccording to one aspect of the present invention, the sketching aid comprises an opaque frame member, having a portion thereof large enough to be readily grasped by a human hand and small enough to be readily supported by a human hand alone, a rigid transparent plate, and means for operatively attaching the plate to the frame so that the plate is held stationary with respect to the frame. Means are also associated with the frame for holding a transparent drawing film, such as an acetate film of substantially the same size as the plate, in face-to-face contact with the plate. Storage means are provided for holding a supply of transparent drawing films, the storage means formed in the portion of the frame adapted to be held by the user's hand. The storage means may comprise means defining a cavity within the frame portion of substantially the same size and shape as a stack of transparent drawing films, the cavity being open on one end thereof to allow sliding movement of the films from the cavity, and a movable abutment mounted adjacent to the cavity for releasably blocking the cavity open end so that the films may not be removed.\nThe means for operatively attaching the plate to the frame may comprise means for pivotally mounted the plate to the frame so that it is movable from a first position wherein the frame portion adapted to be grasped by the user's hand and the plate are in parallel planes, in face-to-face relationship, to a second position wherein the plate makes an angle of between 90.degree. and 180.degree. with respect to the frame portion. A drawing implement may then be held between the frame portion and plate at the hinge, which hinge may comprise a flexible strip.\nThe plate may be considered to have x and y imaginary orthogonal axes disposed in the plane thereof, when held stationary in the frame member, with a z imaginary axis extending orthogonally to the plane of the plate. The means associated with the frame for holding a film in face-to-face contact with the plate comprises means formed intergrally with the frame preventing substantial movement of the film in either direction along the y axis, preventing movement in one direction along the z axis while the plate prevents movement in the other direction along the z axis, and preventing movement along a first direction along the x axis while allowing sliding movement of the film with respect to the frame and plate in a second direction along the x axis, opposite the first direction. A releasable blocking means is disposed at the side of the frame and actuatable by hand to either prevent or allow slidable movement of film in the second direction along the x axis. The dimensions of the film and plate are preferably about 5\".times.8\", with the effective size of the plate and film, as defined by the frame and the attachment means, being about 41/2\".times.65/8\".\nThe device may also be provided with a resting bar means for resting on a surface to steady the frame member during use of the aid. For instance the resting bar may rest on a person's lap, a table, or a like surface, to insure steady support of the aid when being used. Means are provided for mounting the resting bar to the frame member so that the position thereof with respect to the frame member is adjustable, supporting the frame member in each position to which it is moved. Means may be provided supporting a drawing implement on the resting bar means so that when the resting bar means is moved away from the frame member the drawing implement is removable, but when the resting bar means abuts the frame member the drawing implement is not removable from the aid.\nIt is the primary object of the present invention to provide a simple, versatile, complete, and readily transportable sketching aid. This and other object of the invention will become clear from an inspection of the detailed description of the invention, and from the appended claims."}
{"text": "This invention concerns apparatus and methods for measuring the quantity of bubbles present in a liquid containing gas bubbles. The method and apparatus are particularly useful in determining the gas bubble content of liquids turbulently flowing through pipes. The determination is made at a fixed location in the pipe so that the gas bubble content of the flowing liquid can be continuously monitored.\nTechniques for determining the presence or absence of gas bubbles in liquids are known. Apparatus for determining the onset of the formation of bubbles, i.e. cavitation or boiling, but not the proportion of bubbles, has been described by a number of patentees. In U.S. Pat. No. 3,381,525 to Kartluke et al., sound waves are launched into a liquid. The liquid is monitored for sound waves at subharmonic frequencies of the launched sound waves. When subharmonic waves are detected, cavitation is imminent or has begun. In U.S. Pat. No. 3,240,674 to Ledwidge, a similar technique is used. No sound waves are added to the liquid. Instead, the frequencies of sound waves in the liquid are monitored for a selected spectrum peak that indicates localized boiling, a prelude to boiling of the entire liquid volume. U.S. Pat. No. 3,622,958 to Tucker et al. discloses a number of methods of detecting the existence of gas bubbles in a liquid. Waves at a fundamental frequency are launched into a liquid by a first transducer and waves at harmonic frequencies are detected by a second transducer. Detection of harmonic frequency signals indicates the presence of gas bubbles. Alternatively, reflected waves at harmonic frequencies are detected by the same transducer that launches the fundamental frequency wave. In still another embodiment, multiple frequency sound waves are launched into the liquid which is monitored for waves having frequencies equal to a sum or difference of two of the frequencies of the launched waves.\nA complex method of determining the fraction of steam in a steam/water mixture was disclosed by Arave in \"An Ultrasonic Void Fraction Detector Using Compressional Stress Waves in A Wire Helix\" published October 1970 by the Idaho Nuclear Corporation for the U.S. Atomic Energy Commission. In this method, a stress wave is propagated on the surface of a wire helix that is immersed in a liquid-gas mixture. The attenuation of the stress wave from one end of the helix to the other is measured to determine the \"void fraction\", i.e. bubble content, of the liquid in the local volume adjacent the surface of the helix.\nThe known technology does not provide a simple, reliable method of quantitatively measuring the bubble content of a liquid-gas mixture. The invention provides a simple method and apparatus for measuring the bubble content in a liquid and especially in a turbulently flowing liquid."}
{"text": "Recently, discussions for enabling a primarily used frequency band (spectrum) to be used in a second communication service according to the use condition of the frequency band have been in progress. For example, a standard specification for allowing an unused channel (TV white space) included in the frequency band of US digital TV broadcasting to be available for radio communication is under examination in IEEE802.22 working group (see Non-Patent Literature 1 below).\nAccording to a recommendation of the Federal Communications Commission (FCC) on November 2008, the discussions have been directed toward permitting secondary usage of the TV white space by using a communication device that fulfills a predetermined condition and has received an authentication. This recommendation of FCC accepts the standard specification of IEEE802.22, which is the first standardization of secondary usage of the TV white space, and also covers the activities of a new study group in IEEE. Technically, since it is necessary to perform signal detection at a level of −114 [dBm](SNR is about −19 [dB] when Noise Figure (NF) is 11 [dB], for example) using existing technology for example, an auxiliary function such as Geo-location Database Access is expected to be necessary (see Non-Patent Literature 2 below). Also, FCC is under examination for opening a 250-MHz band, which is a part of a 5-GHz band, as a new channel for secondary usage.\nFurthermore, in the EU, there are moves afoot to universally allocate a dedicated control channel referred to as Cognitive Pilot Channel (CPC) for implementing Dynamic Spectrum Access (DSA) under a long-term strategy. Allocation of CPC has been incorporated in the agenda of International Telecommunication Union (ITU)-WP11 in 2011. Technological studies for a secondary usage system that performs DSA are also under way in the IEEE Standards Coordinating Committee (SCC) 41.\nIn general, when a frequency band assigned to a communication service related to primary usage (will be referred to as a first communication service below) is secondarily used, it is important that a communication service related to secondary usage (will be referred to as a second communication service below) not interfere with the first communication service. For this reason, Non-Patent Literature 2 below recommends installation of a data server that receives administrator information, location information, and the like from a secondary usage node which will provide a second communication service and accumulates these pieces of information in a database. In this case, the data server specifies a channel that can be provided for secondary usage according to a request from the secondary usage node, and notifies the secondary usage node of the specified channel. When the secondary usage node uses the channel notified in this way by the data server; interference with a first communication service is prevented."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus equipped with a developer replenishment mechanism that replenishes developer to a development device that develops an electrostatic latent image formed on an image bearing member.\n2. Description of the Related Art\nAn image forming apparatus of an electrophotographic system forms a toner image by developing an electrostatic latent image formed on a photosensitive member using developer (referred to as toner, hereafter) in a development device. Since the amount of toner that can be accumulated in the development device is limited, the toner is replenished to the development device from a container that is detachable to the main body of the image forming apparatus.\nFor example, there is proposed a container that is provided with a rotation portion that is rotated, a pump portion for changing internal pressure of a chamber to discharge toner from the chamber that accommodates the toner, and a converter that converts a rotating motion of the rotation portion into a reciprocating motion of the pump portion (see US 20120014713A1). This container discharges the toner in a chamber by making the pump portion expand and contract according to the rotation of the container. That is, the air that is sucked from a discharging port as the pump portion expands loosens the toner in the chamber, and then, the chamber becomes in a positive pressure state and the air in the chamber forces out the toner, which covers the discharging port, from the discharging port as the pump portion contracts.\nMoreover, there is a known conventional configuration in which a detection device in an image forming apparatus detects a detected portion of a container in order to determine whether the container is mounted to the image forming apparatus.\nThe above-mentioned prior arts suggest a configuration in which the detection device in the image forming apparatus detects the detected portion that rotates with the above-mentioned container in order to check whether the container is mounted to the image forming apparatus and to detect the number of replenishments of toner from the container.\nHowever, when the container rotates just after mounting it to the image forming apparatus so that the detection device is able to detect the detected portion, the toner is discharged from the container whenever the container is mounted to the image forming apparatus.\nAccordingly, when the toner is replenished to the development device from the container so as to supplement consumed toner during an image formation after repeating mounting and dismounting of the container, the amount of toner in the development device increases excessively, which causes an image defect."}
{"text": "The present invention relates to cellular telephone systems. More specifically, the present invention relates to a system for increasing the reliability of the cellular telephone system in environments having substantial multipath propagation or under conditions wherein a large number of mobile telephone units simultaneously attempt to access a base station.\nMany communications systems have multiple transmitters that need to randomly access one or more receivers. A local area network (LAN) is one example of such a multiaccess system. A cellular telephone system is another. In any such system, when several transmitters attempt to transmit simultaneously, the messages may interfere or “collide” with one another. A receiver cannot distinguish among the messages involved in the collision.\nTwo such multiaccess protocols, commonly called the “Aloha” and “Slotted Aloha” protocols, are described in Bertsekas et al., Data Networks, chapter 4, Prentice-Hall, Englewood Cliffs, 1987. In the Aloha protocol, each transmitter may transmit a message at any time. Upon discovering that the transmitted message has collided, the transmitter waits a random delay time and retransmits the message. In Slotted Aloha, all messages fit into a time slot of a predetermined length. Upon discovering that the transmitted message has collided, the transmitter delays a random number of slots and then retransmits the message. In both methods, a random delay is introduced to prevent transmitters from retransmitting simultaneously.\nThe use of code division multiple access (CDMA) modulation is one of several techniques for facilitating communications in which a large number of system users are present. The use of CDMA techniques in a cellular telephone system is disclosed in U.S. Pat. No. 5,056,109 entitled “Method and Apparatus for Controlling Transmission Power in a CDMA Cellular Telephone System” and in U.S. patent application Ser. No. 07/543,496 entitled “System and Method for Generating Signal Waveforms in a CDMA Cellular Telephone System,” now U.S. Pat. No. 5,103,459, both assigned to the assignee of the present invention and incorporated herein by reference.\nIn the above-mentioned patents, a multiple access technique is disclosed where a large number of mobile stations, each having a transceiver, communicate through base stations, also known as cell-sites, using CDMA spread spectrum communication signals. The base stations are connected to a mobile telephone switching office (MTSO), which in turn is connected to the public switched telephone network (PSTN).\nThe use of CDMA spread-spectrum techniques maximizes the number of mobile stations that can communicate simultaneously with the base station because the same frequency band is common to all stations. Each mobile has a pseudonoise (PN) code uniquely associated with it that the mobile station uses to spread its transmitted signal. In the above-referenced patent, this PN code is called the “long PN code.” Once the call has been initiated, i.e., the base station has selected the long PN code corresponding to the transmitting mobile station, the base station can receive and de-spread the signal transmitted by the mobile station. Similarly, the mobile station can receive and de-spread the signal transmitted by the base station. In some systems, the signals may be modulated with a “pilot” PN code as well.\nHowever, for certain types of transmissions, it is advantageous to use a common PN long code, rather than a unique long code for each mobile station. The message transmitted by a mobile station attempting to initiate a call is one example of such a transmission. A mobile station wishing to initiate calls can transmit such requests on a common “access channel” using a corresponding common PN code. The base station can monitor the access channel by despreading the signal using this PN code. The access channel is used because messages such as those for initiating a call are relatively short in comparison to voice transmissions, and a receiver could more easily monitor a relatively few access channels than the large number of unique “traffic channels” with which the mobile stations are associated by their unique PN long codes.\nThe access channel may be used by the mobile station not only to initiate a call, but to transmit any information to the base station at a time other than during a call that has already been initiated. For example, the access channel may be used by the mobile station to respond to an incoming call initiated by a base station over a “paging channel.”\nUnder any of the conditions discussed above, multiple mobile stations may transmit simultaneously on the access channel. When two mobile stations transmit simultaneously and there is no multipath, the transmissions arrive at the base station separated in time by a delay equal to the difference of twice the distance between each mobile station and the base station. Under most operating conditions, it is unlikely that a large number of mobile stations will be at precisely equal distances from the base stations. However, simultaneously transmitted messages would collide if two or more stations are at the same range. Under most conditions, the base station can distinguish among the transmissions because the time between arrivals of the transmissions at the base station exceeds one PN chip.\nSome operating conditions tend to produce collisions. Collisions are likely to occur when a large number of mobile stations approach the edge of a cell simultaneously, a condition causing handoffs of the mobile stations. The access channel transmissions arrive at the base station simultaneously because the mobile stations are at substantially the same distance from the base station when at the edge of the cell.\nIt is also possible that a large number of mobile users would attempt to simultaneously initiate calls for other reasons such as following a natural disaster. The simultaneous transmissions of multiple mobile stations on the access channel may exceed the maximum throughput of the processor in the base station.\nThe probability of access channel collisions increases with an increase in the number of mobile stations and with an increase in multipath reflections. Multipath compounds the problem because, while the main signals of two transmissions may be separated in time by more than one chip, multipath components of the transmissions may not be. Furthermore, as discussed in U.S. Pat. No. 5,109,390, issued Apr. 29, 1992, a base station diversity receiver may have multiple correlators that combine received multipath components to improve message quality. However, ambiguities may exist between multipath components that would reduce the effectiveness of the diversity receiver. These problems and deficiencies are clearly felt in the art and are solved by the present invention in the manner described below."}
{"text": "Orthopedic prostheses are commonly utilized to repair or replace damaged bone and tissue in a human body. A knee prosthesis, for example, can include a tibial component or a femoral component to replace damaged or destroyed bone in the tibia or femur, respectively. Knee prostheses can seek to provide articulation similar to an anatomical articulation of a natural knee joint.\nTotal knee replacement (TKR) surgery can involve the implantation of several components meant to restore the functionality provided by a natural knee. Typical TKR components include a tibial component, a femoral component, and/or an insert or bearing component disposed between the tibial and femoral components.\nIn certain TKR prostheses, which are oftentimes referred to as “posterior stabilized” prostheses, a cam positioned in an intercondylar fossa of a femoral component cooperates with a spine formed in a tibial component to guide or constrain motion within predefined boundaries. Posterior stabilized prostheses can include a spine integrally formed with a tibial bearing insert (sometimes referred to as a meniscal component), which interacts with a cam formed in a femoral component to promote femoral roll back during flexion of the TKR prosthesis. Posterior stabilized prostheses can be appropriate where a posterior cruciate ligament (PCL) is torn or otherwise damaged, or where the PCL is resected during surgery.\nIn addition, some knee prostheses feature a hyperextension stop, which can be a posterior-facing structure formed in a femoral component configured to engage an anterior-facing surface of the tibial spine when the knee prosthesis is in a “full extension” or “hyperextension” orientation. Such a hyperextension stop can also be referred to as an anterior cam. In these knee prostheses, interaction between the femoral component and spine operates as a physical stop against extension of the knee prosthesis past a predetermined level of extension.\nA natural knee experiences internal and external rotation, i.e., rotation about a generally proximal-distal axis, during flexion. Internal/external rotation can be significant in deep flexion where an asymmetric anteroposterior rollback can occur, i.e., the lateral femoral condyle rolls back faster and/or further than the medial condyle. Internal/external rotation is also significant as the natural knee approaches extension (i.e., 0-20 degrees), a phenomenon sometimes referred to as the “screw-home” mechanism. Knee prosthesis designs can seek to accommodate, promote or drive internal/external rotation mimicking the natural knee."}
{"text": "The present invention relates to: a method of manufacturing a semiconductor device; in particular a technology effectively applicable to a method of manufacturing a semiconductor device formed by connecting a wire with a surface of a pad electrode through a conductive layer.\nA technology of adopting a palladium film as a seed metal film and using an etching liquid of an iodine series as the etching liquid is disclosed in Japanese Unexamined Patent Application Publication No. 2010-157683 (Patent Literature 1)."}
{"text": "When a high-precision metal microstructure is mass-produced, a process of lithographie galvanoformung abformung (LIGA) is useful. The LIGA process, which uses synchrotron radiation (SR) ray among X-rays for its high directivity, is characterized in that it allows deep lithography, processing of a structure having a height of several hundreds of micrometers with a precision on the order of micrometers, easy manufacturing of a metal microstructure having high aspect ratio, and the like. Therefore, the LIGA process is expected to find application in various fields.\nThe LIGA process is a processing technique combining lithography, plating such as electroforming, and molding. According to the LIGA process, a resist film is formed on a conductive substrate, for example, and then is irradiated with SR ray through an absorbent mask having a pattern of predetermined shape to perform lithography. Such lithography allows a resist structure (resin mold) to be formed corresponding to the pattern shape of the absorbent mask. A vacant portion of the resist structure is then electroformed to obtain a metal microstructure. Furthermore, a high-precision metal microstructure obtained by further electroforming may be used as a mold to perform molding such as injection molding so as to obtain a micromolded article made of resin. For example, Japanese Patent Laying-Open No. 2002-217461 discloses a method of manufacturing a metal microstructure as a mold by the LIGA process to form a micromolded article.\nHowever, when lithography is used to form a metal microstructure, it can produce only one resist structure (resin mold) at a time. Therefore, lithography is not suitable for mass production. Furthermore, if a resin mold to be electroformed to form a metal microstructure is not fixed firmly with high precision onto a conductive substrate with an adhesive or the like, dimensional change or separation between the resin mold and the conductive substrate occurs in a subsequent step, disadvantageously resulting in a product of lower precision, a product not suitable for commercial use, and the like. Furthermore, in the case where a typical adhesive is applied to the conductive substrate to fix the resin mold thereon so as to solve the problem of a product of lower precision and a defective product, if the adhesive remains at a vacant portion of the layered structure having the resin mold, electromolding cannot be performed or insufficiently performed thereon, disadvantageously resulting in an uneven metal layer formed thereby. Furthermore, if the entire layered structure having the resin mold is etched to remove the adhesive remaining at a vacant portion thereof, the resin mold itself is damaged accordingly to a certain degree.\nAn object of the present invention is to provide a method of manufacturing a metal microstructure by using a resin mold, in which a mild condition can be set so as to cause less damage to a resin mold and a high-precision metal microstructure can be mass-produced by uniform electroforming."}
{"text": "1. Field of the Invention\nThe present invention relates to a device and a method for transposing a matrix of video signals, which may be applied to a video decoder for decoding a matrix of compressed video bitstream such as signals of an MPEG(Moving Picture Expert Group) format.\n2. Discussion of the Related Art\nAs is well known, recently high definition television(HDTV) broadcasting systems have been developed, some of which are put into pilot operation in some of countries. In the pilot operation of the HDTV broadcasting systems, signals according to the HDTV broadcasting systems are transmitted, which are capable of being processed in an HDTV receiver (hereinafter called as HDTV). In the meantime, there are two classes of signals in the HDTV system of the ATSC (Advanced Television Systems Committee) of the U.S.A. One class of signal is an HD class signal for producing images of a high definition, and the other class of signal is an SD (Standard Definition) class signal(hereinafter called as SD signal) for producing images of a standard definition. The SD class signal has a data transmission rate lower than the HD class signal. The HD TV is adapted to process the two classes of video data, basically. Therefore, there can be four cases of signal transmission between broadcasting stations and HD TV\"\"s; the first case is displaying a HD signal as an HD signal, the second case is displaying an HD signal as an SD signal, the third case is displaying an SD signal as HD class signal, and the fourth is displaying an SD signal as an SD signal. Though the first, third and fourth cases have no particular problems in processing the signals with regard to the system performances, the second case does. That is, the second case has problem in that the images appear broken. Therefore, if an HD signal is to be displayed as an SD class signal, it is necessary to degrade performance of a received HD signal to some extent. For example, an HD signal of 20 MHz MPEG format should be altered to an SD class signal of 6 MHz MPEG format. This degradation of the HD signal performance implies a drop in the resolution. In general, the HD signal performance is degraded by removing a portion of frequency band over a certain frequency from the HD signal. An HD signal of an 8xc3x978 matrix having a portion of frequency band over a certain frequency removed therefrom to suit to the SD class signal is thereby altered into an 8xc3x974 HD signal. This 8xc3x974 HD signal can be displayed as an SD class signal, with a resolution degraded compared to the 8xc3x978 HD signal, but without the problem of broken images. In the meantime, an 8xc3x978 matrix SD signal can be displayed on an HDTV as it is. Because of the aforementioned reason, the HD TV should be adapted to process the 8xc3x978 HD signal as well as the 8xc3x974 HD signal. A general HD TV will be explained with reference to FIGS. 1xcx9c3.\nReferring to FIG. 1, upon receiving a compressed video bit stream from a broadcasting station, the general HD TV synchronizes a desired video bitstream in a tuner 10, of the video bitstream received through an antenna Ant, and demodulates in a demodulator 20. This video bitstream, being an MPEG signal, has a form of matrix. Then, the demodulated video bit stream is restored in a video decoder 30 and processed to a displayable condition in a VDP(Video Display Processor) 40. The signal from the VDP 40 is displayed through a displayer 50. In this instance, the VDP 40 processes the decoded signals appropriate to a performance of the displayer 50.\nFIG. 2 illustrates a detail block diagram of the video decoder 30 shown in FIG. 1. Upon reception of a matrix of the demodulated video bit stream from the demodulator 20 in FIG. 1, a VLD(Variable Length Decoder) 21 therein decodes the demodulated video bit stream to provide DCT(Discrete Cosine Transform) coefficients and motion vectors. The DCT coefficients are scanned by an inverse scanner 22 and inverse quantized by the inverse quantizer 23. Then, an IDCT(Inverse Discrete Cosine Transformer) 25 makes an inverse discrete cosine transformation of the inverse quantized DCT coefficients to provide spatial pixel values. Here, before the inverse discrete cosine transformation in the IDCT 25, the inverse quantized DCT coefficients are transposed in a transposer 24 for easy inverse discrete cosine transformation. On the other hand, a motion compensator 27 uses the motion vectors from the VLD 21 in compensating a reference video frame stored in the frame memory in advance, and an adder 28 adds a signal from the motion compensator 27 and a signal from the IDCT 25, to provide an added value to the VDP 40 in FIG. 1.\nIllustrated in FIG. 3 is an example of the transposer 24 for making alternative read/write of two memories 24a and 24b in transposing them.\nReferring to FIG. 3, the transposer 24 is provided with two SRAM(Static Random Access Memory) 24a and 24b each of a 16 bitsxc3x9764 words size and two multiplexed flipflops 24d and 24e. Upon reception of an enable signal in a read/write controller 24c, the read/write controller 24c provides a read/write controlling signal and a selecting signal to the first SRAM 24a and the second SRAM 24b. . The first SRAM 24a and the second SRAM 24b are operative in opposite manner in response to the read/write controlling signal. That is, a first data is written on the first SRAM 24a at the first time, the second SRAM 24b is left at a disabled state in response to the selecting signal. Then, while the first SRAM 24a is read, the second a SRAM 24b is written of a second data. Thus, the first SRAM 24a and the second SRAM 24b are read and written alternatively. The alternative first, second data from the two SRAMs 24a and 24b are transposed by the multiplexed flipflops 24d and 24e. Besides, there are a transposing method disclosed in a U.S Pat. No. 4,769,790 by time delaying and a U.S. Pat. No. 5,418,487 using dual port memories. These two Patents have a disadvantages either in a size of hardware or in a complicated hardware. And, as has been explained, the background art transposer 24 shown in FIG. 2 transposes an mxc3x97m matrix signal(for example, m=8) using two memories of mxc3x97m words capacity each. Because of the use of two memories, such a background art transposer has disadvantages in that the inverse discrete cosine transformer and the video decoder 30 shown in FIG. 2 have increased areas and comparatively low operation speed of 100 M sample/sec.\nAccordingly, the present invention is directed to device and method for transposing a matrix of video signals and a television receiver employing the same that substantially obviates one or more of the problems due to limitations and disadvantages of the related art.\nAn object of the present invention is to provide a method and a device for transposing a matrix of video signals which can reduce a chip area of the inverse discrete cosine transformer and, further, an area of the video decoder.\nAnother object of the present invention is to provide a method and a device for transposing a matrix of video signals which is operable at fast speed.\nOther object of the present invention is to provide a video decoder and a television receiver employing the same method and device.\nAdditional features and advantages of the invention will be set forth in the description which follows, and in part will be apparent from the description, or may be learned by practice of the invention. The objectives and other advantages of the invention will be realized and attained by the structure particularly pointed out in the written description and claims hereof as well as the appended drawings.\nTo achieve these and other advantages and in accordance with the purpose of the present invention, as embodied and broadly described, there is provided a memory part for writing a matrix of video signals thereon, a write control circuit shift and writes rows of the video signals on the memory part by any one unit either of a row unit or a column unit. The video signals written on the memory are read from the memory by one unit different from the unit in the writing either of the row unit and the column unit by a read control circuit for transposing. Rows of video signals received at the next time are shifted, and written on portions of the memory part emptied in sequence due to the shift in the reading by one unit opposite to the unit in the reading either of the row unit and the column unit in sequence. The memory part may includes either one memory on which the matrix of video signals may be written or two memory of less capacity, resulting to a reduction of an occupied area of the memory part. Because rows of video signals received at the next time is shifted, and written on portions of the memory part emptied in sequence in the reading, a speed of the transposing becomes faster.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory and are intended to provide further explanation of the invention as claimed."}
{"text": "1. Field of The Invention\nThis invention relates to a method and apparatus for forming and sealing bags employed for packaging food or other material and, more particularly, to improved method and apparatus for applying an adhesive material to edges to be sealed.\n2. Description of The Prior Art\nIn a typical vertical form and fill bag-making operation for packaging potato chips and similar snacks, apparatus is frequently employed in which film material in sheet or web form is supplied with a continuous or intermittent motion to a forming apparatus which forms the web into a tubular shape with either overlapping or fin-type edges. Glue or other adhesive is often pre-applied to one or more edges of the film. The film, after being formed into the tubular shape, is passed through a sealing apparatus in which heat and pressure are applied to seal the edges together. In overlapping edges the adhesive is applied to edges on opposite faces of the web; in fin-type seals the adhesive is applied to edges on the same face of the web.\nIn one form of such apparatus, a transverse seal is made to close the top of one bag and form the bottom of the succeeding bag. A downward pull is applied to the transversely sealed area and the overlapping edges of the tubular film are drawn along a device to heat-seal the longitudinal edges of the bag. The bag is filled and a transverse seal is formed at the top of the bag. The formed bag is severed, and the process is then repeated by the apparatus.\nThe film material employed should be one which adequately protects the food material or other product being packaged, is pleasing in appearance, is strong enough to withstand shipping and handling environments, and is amenable to forming and sealing, among other characteristics. The film material often used in the apparatus for forming bags may include a laminate having for example as one component a layer of polyvinylidene chloride class (e.g., that material sold commercially under the trademark \"Saran\") and as the other component a layer of uncoated polypropylene. For a lap seal, i.e., where one edge of the sheet material overlaps the other, this results in a \"Saran\" surface facing the uncoated polypropylene surface in the area of the overlap. These materials are not compatible for heat sealing to one another. Consequently, where heat sealing is involved, it is necessary to employ a separate material secured to one or more of the overlapping edges to insure that the surfaces facing one another for bonding have similar or compatible characteristics.\nWhere a thermal, or heat, seal is desired using dissimilar materials a separate strip of thermal seal material is bonded along one edge of the web by the manufacturer of the film to assure thermal bonding compatibility. This separate strip is oriented in the machine direction along the edge of the web so that as the forming apparatus creates the tube the strip applied, for example, to outer face of the polypropylene will provide a more compatible material with respect to thermal bonding capability facing the \"Saran\" layer side of the web in the vicinity of the overlap. This facilitates forming a thermal seal when the edges are pressed together and heated.\nIn addition to the time and expense required in applying this separate strip, the thickness of the strip causes a number of other problems in handling and forming of the sheet material into filled bags. Buckling can occur in the vicinity of the web material where the strip has been applied, causing the material to pucker when sealed. When the web material is rolled, one end of the roll, where the strip is applied, tends to be thicker than the other, causing the web to stretch. As a result uniformity and reliability in forming and sealing the bags is lost or at least substantially impaired. There is also a problem with air bubbles getting into adhesive seals."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to apparatuses for securing loads. In particular, the present invention relates to buckle securing methods, apparatuses, and systems.\n2. Description of the Related Art\nLarge and heavy pieces of cargo are often made up of smaller components that must be fixed together before transport. For example, lumber, long planks, sheets of plywood, and large beams or columns are typically strapped together before being placed on a land transporter (e.g., trucks), a water transporter (e.g., ships), and/or an aerial transporter (e.g., planes or helicopters) for shipping.\nThe straps used are typically made of relatively expensive webbing (e.g., a polyethylene webbing) that comes in lengths of about 10 meters.\nTypical couple/buckle devices/systems are tied/secured to one end of the webbing. After the webbing is wrapped around the objects (for strapping them together), the other end of the webbing is attached to the buckle. Any excess webbing (i.e., the tag end) is cut off and discarded (i.e., wasted).\nIn addition, an installer needs to position the typical buckle in such a location that the installer has enough leverage to pull and tighten the webbing around the objects. If the webbing is not cinched tight enough the objects will not remain secured by the prior art buckle and webbing. Objects that are not properly secured can become separated during transport, separated when deployed from the transport, lost, and/or damaged. When not cinched tight enough the installer will need to loosen and reposition the buckle so that the new position allows the installer enough leverage to appropriately tighten the buckle/webbing system.\nFurther, it is difficult to cinch less expensive webbing (e.g., polypropylene webbing) enough (using prior art buckles) to secure objects.\nAfter the prior art buckle systems are deployed, troops typically cut and discard the webbing and buckle.\nThus there is a need in the art for a buckle system that is easy to install, is less likely to need repositioning for leverage, allows less waste of webbing material, is easily reusable, and can be used with less expensive webbing."}
{"text": "Infectious complications of trauma injuries often require treatment of multi-drug resistant bacterial pathogens, and their secreted toxins. Bacterial toxins play a significant role in pathogenesis, virulence, tissue damage and inflammation. War wound infections have long posed a major challenge for military medicine and, infectious complications will remain a major cause of morbidity and mortality.\nUS casualties from Iraq and Afghanistan often have wounds that are colonized or infected with multidrug-resistant strains of Staphylococcus aureus. Other frequently identified resistant strains of bacteria are Klebsiella pneumoniae, Pseudomonas aeruginosa and Acinetobacter baumannii. Overuse of broad-spectrum antibiotics may be the primary cause in emergence of these resistant strains [1]. In addition, these pathogens secrete toxins which enhance virulence and tissue damage. Such factors likely lead to hard to heal wounds.\nWhen toxin producing pathogens invade and multiply in tissues, it is likely that during infection, sensitive cells will be exposed to their products, such as toxins that provoke active cellular responses. Synthesis and secretion of inflammatory cytokines is, of course, one of the countless possible reactions that will follow as a consequence of infections. S. aureus expresses numerous virulence factors including α-toxin/hemolysin that are aimed at weakening the immune system by direct inactivation of key innate response molecules. As a pore forming factor, α-toxin plays a major role in the pathogenesis of S. aureus infections and, therefore, is an important target for therapeutic intervention [8, 11]. Although early intervention with antibiotics can reduce the bacterial burden, inhibitors of secreted toxins could offer a therapeutic advantage in the later stages of infection involving toxemia. Currently, no such inhibitors are available to neutralize α-toxin mediated cytolytic and cytotoxic effects. Therefore, the present invention was undertaken to identify DNA based inhibitors/aptamers to detoxify α-toxin induced toxicity in S. aureus infections.\nAptamers represent a class of macromolecules with attractive drug properties which could be useful in the treatment of a variety of ailments. Aptamers are nucleic acids having specific binding affinity to multiple types of molecules through interactions other than classic Watson-Crick base pairing. Aptamers, like peptides generated by phage display or monoclonal antibodies (“MAbs”), are capable of specifically binding to selected targets, and modulating the target's activity. For example, binding of aptamers to a target may block the target's ability to function. Created by an in vitro selection process from pools of random sequence oligonucleotides, aptamers have been generated for over 100 proteins including growth factors, transcription factors, enzymes, immunoglobulins, and receptors. A typical aptamer is 10-15 kDa in size, consists of 30-45 nucleotides, binds its target with sub-nanomolar affinity, and discriminates against closely related targets. Aptamers will typically not bind other proteins from the same gene family. A series of structural studies have shown that aptamers are capable of using the same types of binding interactions, such as hydrogen bonding, electrostatic complementarity, hydrophobic contacts and steric exclusion that drive affinity and specificity in antibody-antigen complexes.\nAptamers have a number of desirable characteristics for use as therapeutics and diagnostics including high specificity and affinity, biological efficacy, and excellent pharmacokinetic properties. In addition, they offer specific competitive advantages over antibodies and other protein biologics, for example speed and control. Aptamers are produced by an entirely in vitro process, allowing for the rapid generation of initial leads, including therapeutic leads. In vitro selection allows the specificity and affinity of the aptamer to be tightly controlled, and allows the generation of leads, including leads against both toxic and non-immunogenic targets.\nAptamers as a class have demonstrated little or no toxicity or immunogenicity. In chronic dosing of rats or woodchucks with high levels of aptamer (10 mg/kg daily for 90 days), no toxicity is observed by any clinical, cellular, or biochemical measure. Whereas the efficacy of many monoclonal antibodies can be severely limited by immune response to antibodies themselves, it is extremely difficult to elicit antibodies to aptamers most likely because aptamers cannot be presented by T-cells via the MHC and the immune response is generally trained not to recognize nucleic acid fragments.\nAptamers (protein-binding oligonucleotides) have great potential as a new class of targeted therapeutics. These molecules are artificial, single-stranded nucleic acid ligands (DNA and RNA) generated against amino acids, proteins, drugs, viruses and whole cells [2]. They are discovered by an in vitro directed evolution process called SELEX but are ultimately chemically synthesized like small molecule therapeutics using solid-phase techniques that are easily scaled up and allow for modification and conjugation strategies [3, 4, and 5]. Aptamers emerged as a new class of molecules that rival antibodies in both therapeutic and diagnostic applications. One of the advantages of DNA-based aptamer drugs over currently available low molecular weight pharmaceuticals or antibodies is their selective recognition of molecular targets which imparts tremendous specificity of action with binding affinities in the low nano- to picomolar range. Because of their unique three dimensional structures, aptamers are able to bind to functional domains of the target, thereby modulating the biological function of the molecule. Their small size enables them to access epitopes that might otherwise be blocked or hidden. Aptamer selection and identification is much faster than monoclonal antibodies and allows optimization of binding affinity by successive rounds of selection and screening. Antibodies generally elicit strong immune responses due to their size, whereas aptamers show little or no immunogenicity. Further, aptamers are thermally stable and preserve their structure over a wide range of temperatures, exhibit low toxicity, and can be generated to recognize a wide range of targets, including small molecules and ions.\nIn spite of the infancy of aptamer therapeutics, these molecules have produced striking results in both preclinical and clinical studies. Several aptamers are currently in clinical trials, or being evaluated, for different diseases [25, 26]. To date, there are a few aptamers identified against bacterial and plant toxins that have the potential to inhibit the biological functions of the respective toxin both in vitro and in vivo [27-29]. Recent work by Liang et al. further reinforces the therapeutic potential of aptamers targeting rabies virus and demonstrated the aptamer mediated inhibition of rabies viral multiplication both in vitro and in vivo [30]. Most notably, the only DNA aptamer identified against one of the S. aureus exoproteins was against enterotoxin B (SEB) and is designated as APTSEB. This aptamer was developed not as a therapeutic, but to be used in the diagnostic platform to detect SEB [31].\nThe present invention provides novel small drug molecules or aptamers and method of using them to neutralize the lethal effects of Staphylococcus aureus alpha toxin for treatment of S. aureus infections."}
{"text": "The invention relates to a stacking disk oil cooler for a motor vehicle engine, comprising several trough-shaped plates which, with their upright edges situated inside one another, are stacked in a spaced manner and are soldered together for forming adjacent hollow chambers, the lowest plate being mounted on a base plate.\nStacking disk oil coolers of this type are known from European Patent Document EP 0 623 798 A2. In the case of constructions of this type, the lowest of the trough-shaped plates, which is mounted on the base plate, forms the weakest link with respect to the stability. Oil coolers are generally fastened directly by means of the base plate on the engine block and, during the operation of the engine, are therefore subjected to extremely high lateral acceleration forces. As a rule, stacking disk oil coolers of the above-mentioned type do not withstand such forces. During the operation, the lowest disk may deform so that the effect of the cooler becomes doubtful.\nIt is an object of the present invention to further develop a stacking disk oil cooler of the initially mentioned type such that it corresponds to the stability requirements also when it is mounted directly on the engine block so that it will not be necessary to use other cooler constructions for these types of mountings.\nFor achieving this object, it is provided according to the invention in the case of a stacking disk oil cooler of the initially mentioned type that a reinforcing disk is arranged between the base plate and the lowest plate, which reinforcing disk is provided with a rim surrounding the edge of the lowest plate. By means of this arrangement, the thin edge of the lowest plate, which when installed directly on the engine block is subjected to the risk of buckling, is reinforced so that in a relatively simple manner stacking disk oil coolers are also suitable for a direct mounting on the engine block. Because of their higher stability, oil coolers according to the invention can also be used for different purposes if higher stability is important, according to other contemplated embodiments.\nAs a further development of certain preferred embodiments of the invention, the rim has the shape of upright tabs which rest only against the longitudinal and transverse sides of the edge of the rectangular plate. This permits a very simple development of the reinforcing disk which, however, nevertheless results in the desired reinforcing effect. In a further development of certain preferred embodiments of the invention, the tabs may project from the reinforcing plate at an angle which is adapted to the angle of slope of the edges of the plate, and are correspondingly bent for this purpose. In this case, it is not necessary that a contact takes place in the bending area with the edge bend of the lowest plate. It is sufficient for the tabs to rest against the longitudinal and transverse sides of the upright edge of the lowest plate so that, as a result, they project into an area in which, because of its construction, the stacking disk oil cooler already has a higher stability. An adaptation of the tabs of the reinforcing disk also in the bending area, which is more difficult to implement, is usually unnecessary but, if the use is extreme, can also be implemented.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "Implantable pacing devices often include electrodes to deliver pulses and/or to sense cardiac activity in response to delivered pulses. Cardiac activity related to a delivered pulse is typically known as an “evoked response” (e.g., an electrical signal arising from atrial or ventricular cardiac tissue depolarization in response to delivery of a pacing pulse). However, many issues complicate sensing and/or detection of evoked responses. For example, post-pulse electrode polarization can interfere with detection of an evoked response and/or produce a “polarization artifact” in a detected signal. Post-pulse electrode polarization results primarily from capacitive charging of an electrode-electrolyte interface during delivery of a pacing pulse. Upon termination of the pacing pulse, the post-pulse electrode polarization decays over time, generally in an exponential fashion like a capacitor. Characteristics of post-pulse electrode polarization generally depend on a variety of parameters, such as, electrode materials, electrode geometry, tissue characteristics, tissue contact, stimulation energy, and others, many of which vary over time. Consequently, an elaborate characterization of post-pulse electrode polarization is impractical, especially when one considers resource limitations inherent in implantable pacing devices.\nU.S. Pat. No. 6,163,724, entitled “Microprocessor capture detection circuit and method”, to Hemming, et al. ('724 patent), addresses post-pulse electrode polarization through use of filtering. More specifically, the '724 patent discloses “an adaptive nonlinear filtering technique referred to as ‘Negative Peak Tracking’ (or ‘NPT’) that removes the initial residual [post-pulse] polarization signal, and then passes only that portion of the sensed signal where a change in the sign of the slope occurs” (col. 7, lines 4-8). According to the '724 patent, such filtering, in combination with “[j]udicious selection of comparator threshold levels by a user[,] improves the reliability of event discrimination [capture versus non-capture events]” (col. 7, lines 26-28).\nAnother issue in detection of evoked responses stems from differences in ventricular and atrial pacing. Thus, approaches to detection of evoked responses in ventricular pacing may not apply directly to detection of evoked responses in atrial pacing. For example, U.S. Pat. No. 5,713,934, entitled “Evoked and spontaneous cardiac activity detection in a dual-chamber electronic pacemaker and method”, to Leckrone, ('934 patent) discloses a pacing system that uses ventricular pacing electrodes and atrial pacing electrodes for sensing. According to the '934 patent, a ventricular evoked response typically peaks at about 30 milliseconds following a ventricular pulse and ranges in amplitude from about 3 millivolts to about 20 millivolts; whereas, an atrial evoked response typically occurs within 20 milliseconds of an atrial pulse with an amplitude generally smaller than that of a ventricular evoked response.\nWhile other pacing and sensing electrode configurations are possible, none of the patents referred to herein discuss the relationship between post-pulse timing of an atrial evoked response and electrode configuration. For example, U.S. Pat. No. 4,549,548, entitled “Pacemaker system with automatic event-programmed switching between unipolar and bipolar operation”, to Wittkampf, et al., issued Oct. 29, 1985 ('548 patent), discloses a pacemaker system capable of both unipolar and bipolar sensing wherein “considerations for unipolar and bipolar sensing vary at different times in the pacing cycle, dependent upon the next anticipated event” (col. 2, II. 12-19). Thus, the '548 patent focuses on the relationship between electrode configuration and “the next anticipated event” and not other considerations such as, but not limited to, electrode polarization. Likewise, U.S. Pat. No. 4,858,610, entitled “Detection of Cardiac Evoked Potentials”, to Callaghan, et al., issued Aug. 22, 1989 ('610 patent), fails to discuss the relationship between post-pulse timing of an atrial evoked response and electrode configuration.\nU.S. Pat. No. 5,873,898, entitled “Capture Detection Circuit for Pulses and Physiological Signals”, to Hemming, et al., issued Feb. 23, 1999 ('898 patent), discloses a system for pacing and sensing and presents data from a canine ventricular pacing study for pacing and sensing in both unipolar and bipolar electrode configurations. More specifically, in Table 3 of the '898 patent, “Event 15” uses a ring-to-can electrode configuration for ventricular sensing (col. 25, II. 56-57). While the '898 patent alludes to enhanced accuracy for tip-to-can (unipolar) ventricular sensing when compared to tip-to-ring (bipolar) ventricular sensing (col. 29, II. 51-55), the '898 patent makes no further mention of “Event 15”. As shown in Table 4 of the '898 patent, for Event 15, the capture detection circuit had a success rate index of less than one for three of the six pacing settings tested (col. 28, II. 36-40). Therefore, the '898 patent suggests that ring-to-can sensing is not beneficial for detection of ventricular evoked responses.\nNone of the aforementioned patents discuss an atrial pacing and sensing system that diminishes and/or eliminates the effects of electrode polarization through electrode configuration. In particular, none of the aforementioned patents disclose an atrial pacing and sensing system wherein detection of an atrial evoked response occurs more than approximately 20 milliseconds after administration of an atrial pulse and/or wherein an atrial evoked response is relatively independent of pulse power."}
{"text": "Multiple computer databases or databases with multiple partitions pose special problems for database designers. In many cases, it is desirable to have a set of databases that hold all data, but where each user of the system has a unique view into that data. Stated another way, each user would have access to just a subset of all available data. However, database designers are limited by software applications that each require the appearance of a single database.\nIn computer networking environments which store network topology information, one problem is to maintain multiple, independent databases where each single network topology database is located on a physically distinct machine. Topology processes on each machine typically write to their own database and are not aware of each other. To accommodate current software applications, the databases must have the appearance of a single database or single datascape to maintain backward compatibility and eliminate the need for special applications to handle multiple databases explicitly.\nWhen displaying network information to system users, a union of all data in view must be formed across the multiple databases, resolving duplicated and/or missing data. To give the appearance that the data came from a single, self-consistent database, the resulting datascape typically must be uniform, logically consistent, and homogeneous. In addition, information must be gathered from multiple databases in an efficient and seamless manner.\nTherefore, there is a need for a way to implement multiple computer databases or databases with multiple partitions where a single, unique view is presented to each user of the system. There is a further need to have the database appear as a single database to software applications. There is a further need to gather information from databases in an efficient and seamless manner."}
{"text": "The precise spatial reposition of the brain to millimeter accuracy is necessary for implementing certain therapeutic techniques, for example surgery in stereotaxic conditions, or repeated irradiation for tumour-related illnesses or vascular malformations (angiomas).\nPositioning of a reference frame, which may or may not be metal, fixed to the neurocranium allows investigations to be performed in a topographically superposable manner: arteriography, ventriculography, body scanning, MRI, and performing treatments, which constitutes the principle of stereotaxy. But this assumes that the frame, and therefore the fastening of the cranial anchoring, is repositioned each time.\nTwo technical methods are possible for this purpose:\neither a position-finding process which is redone each time, using three-dimensional coordinate computation. This computation is currently relatively easy but the refitting does not always have millimeter precision, taking into account that the needles of the frame are difficult to reposition in the holes of the bone. PA1 Or precise repositioning, to millimeter accuracy, of the frame in the osseous structure of the neurocranium using percutaneous bone screws.\nThis has obvious resulting drawbacks from the point of view of the patient: his neurocranium must actually be subjected to piercings, incisions of the scalp and local anaesthetics repeated on each refitting, giving rise to infection risks and painful discomfort. For the surgeon, repeating these fitting operations requires considerable time trying to find the position accurately.\nThe latter technique, undoubtedly the most accurate, demands iterative refitting of the frame on the neurocranium, on percutaneous screws (that is to say, screws whose head projects out of the skin of the patient) which are anchored in the osseous structure, for several days on end or possibly several days per week for several weeks, during sequential irradiations."}
{"text": "Security systems have traditionally been utilized for detecting, removing, etc. unwanted data (e.g. malware, spyware, unsolicited data, etc.). Oftentimes, it is advantageous for separate security systems to communicate, such as for sharing information regarding known, detected, etc. unwanted data. Unfortunately, traditional security systems have exhibited various limitations with respect to communication capabilities thereof.\nJust by way of example, traditional security systems are generally periodically updated and upgraded, and the frequency of such updates and upgrades is increasing such that new types of systems are frequently becoming available. While new security systems may sometimes be capable of being programmed to communicate with old security systems via protocols existing in the old security systems, requiring such programming is restrictive since the protocols utilized must be defined at the time when the old security system was developed. In addition, when the old security systems need a way to communicate with newer security systems, communication therebetween is impossible since the old security systems are not already equipped with the communication protocols used by the newer security systems.\nA similar problem exists for two security systems developed independently (e.g. as a result of different developers or, simply, due to a lack of internal communication between such developers). Both security systems may be released, but to pass messages from one another both products may be required to be re-released, which may be very costly. There is thus a need for addressing these and/or other issues associated with the prior art."}
{"text": "Living and non-living catalysts have been used to polymerize olefins. In living polymerization, each catalyst molecule initiates a growing polymer chain that does not undergo chain transfer or termination reactions while monomer is present. One can determine whether or not living polymerization has occurred by comparing the number of initiator molecules with the number of polymer chains produced in the final polymer. These two numbers should be equivalent for true living polymerization. If there are a substantially greater number of polymer chains as compared to initiator molecules, then the polymerization is not living.\nThe basic components of known living polymerization systems include a Lewis acid, a tertiary alkyl initiator molecule containing a halogen, ester, ether acid, or alcohol group, and an electron donor molecule such as ethyl acetate.\nLewis acids of known living polymerization systems include titanium tetrachloride (TiCl.sub.4), boron trichloride (BCL.sub.3), tin tetrachloride (SnCl.sub.4), iron trichloride (FeCl.sub.3), aluminum trichloride (AlCl.sub.3) and the like. These compounds have been described in U.S. Pat. Nos. 4,910,321 and 4,929,683 and European patent application 341,012 for use in living polymerization of olefins. The exact combination of these elements, however, varies with each system. The tertiary alkyl initiators typically used in living polymerization systems are represented by the formula: ##STR1## wherein R.sub.1, R.sub.2 and R.sub.3 are a variety of alkyl or aromatic groups or combinations thereof, n is the number of initiator molecules, and X is the functional group upon which the Lewis acid affects a change to bring about the carbocationic initiating site. The functional group is typically a halogen, ester, ether, alcohol, or acid group depending on the Lewis acid employed.\nAs discussed in U.S. Pat. No. 5,169,914, the chosen electron pair donor component of the above systems is believed to directly relate to the ability of these catalysts to stabilize the carbocation formed and to generate living conditions. The electron donor number concept has been used to explain the activity of catalyst systems which employ ether and ester initiators. It is believed that the formation of in situ electron pair donors are responsible for the catalyst characteristics. However, the role of the electron donor is still uncertain and has been challenged. See, M. Gyor, L. Balogh, H. C. Wang, R. Faust, Polym. Prepr. Amer. Chem. Soc., 33(1), 158 (1992).\nSuch living polymerization systems are useful for production of narrow molecular weight distribution polymers, however, these systems are very sensitive to certain impurities. For example, the presence of water in Lewis acid catalyzed systems has been considered detrimental because water is associated with chain transfer dominated, nonliving polymerization which results in broad molecular weight distribution polymers i.e., polymer with an Mw/Mn greater than 2.0. There is a decrease in desirable physical properties when the polymer has a broad molecular weight distribution. Recently we have found that water can be used to initiate the living polymerization of isobutylene provided that the correct Lewis acid to water ratio is used. Additionally, in co-pending U.S. application Ser. No. 08/044,862 the addition of an amine to this process was found to modify the overall amount of Lewis acid required to bring about living conditions.\nCatalyst systems based on BCl.sub.3 and TiCl.sub.4 using a number of combinations of the previously described components have very similar process characteristics. First, the Lewis acid concentrations must exceed the concentration of initiator sites by 16 to 40 times in order to achieve 100 percent conversion in 30 minutes (based upon a degree of polymerization equal to 890) at -75.degree. to -804.degree. C. Much longer polymerization times are required for higher polymerization temperatures. These catalyst systems are also typically used with solvents such as methyl chloride.\nFor an industrially applicable process, these catalysts and polymerization conditions fall short of commercial usefulness. Improvements in these systems would include reduction in polymerization time while maintaining narrow molecular weight distribution."}
{"text": "This invention relates to a semiconductor integrated circuit device. More particularly, the present invention relates to a technique which will be effective when applied to a semiconductor integrated circuit device having self-alignment type bipolar transistors inside a region defined by an element isolation region for isolating semiconductor elements from each other.\nThe bipolar transistor which is now being developed by the inventor of the present invention employs an SICOS (Side Wall Base Contact Structure) structure disclosed in Japanese Patent Laid-Open No. 95662/1988. A so-called \"self-alignment type\" bipolar transistor employing this SICOS structure has a low parasitic capacitance formed in each of the emitter, base and collector regions and can therefore obtain high frequency characteristics at a low current.\nThe self-alignment type bipolar transistor is formed in a protruding (convex) island region which is in turn formed in an active region of a semiconductor substrate. The operation regions of the emitter, base and collector regions are formed vertically in the protruding island region from its main plane in the direction of depth of the semiconductor substrate. A graft base region is disposed on the side wall (shoulder portion) of the protruding island region and a base lead-out electrode is connected to this graft base region in self-alignment therewith. An emitter lead-out electrode is connected to the emitter region. Each of the emitter region and emitter lead-out electrode is formed in self-alignment with the base lead-out electrode.\nThe island region in which the self-alignment type bipolar transistor is to be formed is encompassed by the element isolation region and is electrically isolated from other island regions disposed on the semiconductor substrate. The element isolation region consists primarily of thin trenches, an insulating film disposed inside and along each thin trench and a buried material buried into the thin trench through the insulating film. The thin trench is formed at a position spaced-apart by a predetermined distance from the protruding island region. This thin trench is formed by etching anisotropically the semiconductor substrate by RIE or the like. The insulating film formed along the surface of the thin trench is a silicon dioxide film formed by thermally oxidizing the semiconductor substrate on the surface of the thin trench. The buried material is formed selectively only inside the thin trench by depositing a polycrystalline silicon film on the entire surface of the semiconductor substrate by CVD (Chemical Vapor Deposition) and etching back the polycrystalline silicon film thus deposited. The device isolation region of this kind has a small occupying area because the thin trench is formed by anisotropic etching having a small side etching quantity, and can therefore improve the integration density of the semiconductor device."}
{"text": "The present invention is directed toward optical communication systems incorporating dispersion compensation and particularly such systems having an adjustable amount of dispersion compensation and to components of such systems.\nOptical signals transmitted in a fiber optic communication system typically constitute a series of pulses of digital information. Although the pulses are usually at a single nominal wavelength, each pulse is actually composed of different spectral components. The spectral components of each pulse propagate through the transmission fiber at different speeds with higher frequency components traveling slower than lower frequency components in non-dispersion shifted optical fiber. This effect, known as xe2x80x9cchromatic dispersionxe2x80x9d, can result in spectral components of one pulse arriving at a receiver at substantially the same time as a succeeding pulse, thereby causing degraded receiver sensitivity. Chromatic dispersion becomes increasingly pronounced at higher bit rates, e.g. those associated with synchronous optical network (SONET) OC-192 transmission speeds.\nTypically, optical signals propagating through transmission optical fiber experience a positive dispersion. Accordingly, dispersion compensated fiber (DCF) having a negative dispersion can be coupled to the transmission optical fiber in order to offset the chromatic dispersion thereof. The amount of dispersion experienced by an optical signal depends on the distance traveled through the transmission optical fiber. Shorter lengths of transmission optical fiber require less compensation than longer lengths. Likewise, longer lengths of DCF are used when more compensation is required, and shorter lengths of DCF are provided when less compensation is needed. Each segment of transmission optical fiber, however, requires a unique length of DCF in order to provide an appropriate amount of dispersion compensation.\nIn a fiber optic network, many different lengths or spans of transmission optical fiber are used to connect various points within the network. According to the conventional approach to dispersion compensation described above, a unique length of DCF must be cut for each span. If the network contains many spans, a corresponding number of DCF fibers must be provided, each having a different length. As a result, the conventional approach to dispersion compensation in such instances can be inefficient and time consuming.\nDispersion compensating modules have been proposed including segments of DCF fiber of varying length to provide tailored dispersion compensation. Some of the DCF segments have a positive dispersion, while others have a negative dispersion. Selected DCF segments are coupled to one another to provide a desired net dispersion to offset the dispersion associated with the transmission optical fiber. Thus, rather than provide a unique segment of DCF for each span, the same dispersion compensationmodule can be used for spans of varying length and fiber types by simply connecting appropriate segments of DCF within the module.\nThe segments of DCF to be included in the dispersion compensation module are typically wound about standard-sized spools. The spools, in turn, are housed within a package. If used in the United States, however, the package should conform to a Network Equipmentxe2x80x94Building System (xe2x80x9cNEBSxe2x80x9d) standard, as well as other standards dictating fiber packaging dimensions in order to be a commercially attractive. Moreover, the spools should be readily coupled to one another in order to provide the desired compensation, but still conform to these, as well as any other pertinent, standards.\nConsistent with the present invention, a dispersion compensation module is provided which comprises a housing accommodating one or more spools of DCF. Each spool is provided in a respective package, which can be slidably inserted into or removed from the housing along a track provided within the housing. xe2x80x9cJumpersxe2x80x9d or relatively short pieces of fiber can be used to interconnect the spools through adapters external to the packages. As a result, spools of DCF can be readily coupled to one another to provide a desired amount of dispersion compensation. Moreover, the adapters are geometrically configured so that connections can be made to the spools without excessively bending fiber interconnecting the DCF and the DCF itself, while permitting the module to conform to fiber optic packaging standards."}
{"text": "Natural language processing (NLP) is a field of computer science, artificial intelligence, and linguistics concerned with the interactions between computers and human (natural) languages. As such, NLP is related to the area of human-computer interaction. Many challenges in NLP involve natural language understanding that would enable a computer to derive meaning from human or natural language input."}
{"text": "Imaging X-ray projections can be acquired via a multiplicity of medical devices for X-ray imaging, for example computed tomography systems, C-arm angiography systems, radiography devices or mammography devices. The medical devices used for X-ray imaging include at least one X-ray source. The X-ray source illuminates an examination subject and the X-ray detector positioned behind it. The X-ray beams are partially absorbed or attenuated by the examination subject. The X-ray detector is able to detect X-ray beams that pass through the examination subject.\nIn current-generation computed tomography systems, X-ray detectors having a scintillator material are used for converting X-ray radiation into optical signals. The optical signals are converted into electrical signals via photodiodes. The spatial resolution is limited to 15 to 20 line pairs per cm. In addition to the constraints in relation to image quality, the possibilities for deployment in clinical applications are also limited by the restricted spatial resolution.\nAn improvement in image quality, in particular an improved spatial resolution, would be desirable in the imaging of small structures, for example stents in the coronary arteries, in the assessment of in-stent stenoses or in imaging during interventions using needles for example. Using needles during an intervention can lead to partial volume artifacts being generated due to the inadequate spatial resolution. The partial volume artifacts are manifested as dark streak artifacts at the tip of the needle, thereby making an accurate position determination more difficult. In present-day computed tomography systems, small structures, such as, say, stents in the coronary arteries, can be examined only up to a minimum size. The partial volume artifacts produced when needles are used during interventions can only be reduced or suppressed to an inadequate extent via software-based image corrections."}
{"text": "1. Field of the Invention\nThe present invention relates to an improved process for the conversion of an aromatic hydrocarbon in the presence of a titanium tetrachloride catalyst system.\nThe invention will be described with reference to alkylation, i.e., the synthesis of cumene by alkylation of benzene with propylene in the presence of the catalyst, but will find use in alkylaromatic transalkylation and isomerization as well.\n2. Description of the Prior Art\nConversion of aromatic hydrocarbons is well known in industry. Some of the aromatic conversion reactions which occur include alkylation of aromatic hydrocarbons with an alkylating agent such as an olefin, disproportionation of alkylaromatic hydrocarbons and isomerization of alkylaromatic hydrocarbons such as xylenes.\nOf special interest, has been the propylation of benzene to cumene. Cumene is used for the production of phenol and acetone. Cumene is also dehydrogenated to form methylstyrene, in a process similar to that used to convert ethylbenzene to styrene. Cumene is also used as a blending component in aviation gasoline because of its high octane number. The consumption of cumene in the U.S.A. was about 350,000 metric tons in 1968. Of this total, 94% was used for the production of phenol or acetone.\nIt is well known that cumene can be synthesized from benzene and propylene using a catalyst of AlCl.sub.3, SPA or BF.sub.3. SOA is a generally accepted abreviation for solid phosphoric acid catalyst, or phosphoric acid which is adsorbed on kieselguhr or other support.\nAlCl.sub.3 is a very popular alkylation catalyst, because of its high activity. Unfortunately, the catalyst operates as a slurry or sludge which is messy to handle on a commercial scale, and also is corrosive. The highly reactive nature of this Friedel-Crafts metal halide catalyst, AlCl.sub.3, is desirable when attempting to alkylate benzene with ethylene, because less active catalyst systems do not work. However, for alkylation with propylene such highly reactive systems are not necessary.\nThe AlCl.sub.3 catalyst, although not equivalent to a Ziegler-Natta catalyst, is similar, at least for some olefin reactions. AlCl.sub.3 will promote polymerization of both ethylene and propylene, and alkylation of benzene with both ethylene and propylene. These catalysts also promote transalkylation of alkyl groups. This is in contrast to slightly less active systems, such as SPA catalyst, which catalyzes alkylation of benzene with propylene, but does not catalyze transalkylation satisfactorily.\nAnother highly selective catalyst system has been developed for the alkylation of benzene with olefins. This catalyst comprises boron trifluoride. The boron trifluoride catalyst system is exceptionally active and permits operation with dilute olefin streams, but it requires the continuous addition of BF.sub.3 to maintain catalyst activity. High catalyst activity also leads to oligomerization of olefins, so contact time of olefins with BF.sub.3 catalyst should be as short as possible. This catalyst is also exceptionally water sensitive, as water not only destroys the catalyst, but produces very corrosive solutions which attack downstream processing units. BF.sub.3 also frequently appears in the product, and must be removed therefrom.\nBecause of the interest in alkylating benzene with olefins, and because of the inadequacies of existing catalyst systems, I studied the work that others had done, and made exhaustive investigations to determine if it would be possible to find a catalyst which would have the activity and selectivity required to produce an acceptable cumene product, while making maximum use of existing petroleum resources.\nA highly active catalyst was sought, to permit operation at lower temperatures with less utility cost, cost of construction, and to operate with less catalyst. In new units this would mean smaller, and cheaper reactor vessels, while in existing units it would mean that an increase in capacity could be obtained by changing catalyst in an existing reactor vessel, with minor modifications.\nHigh selectivity is necessary, not only to permit operation with feedstreams which are not 100% pure olefin, but also to maximize production of the desired product, and to minimize production of polymerized olefins, or polyalkylated aromatic compounds.\nAccordingly, many catalyst systems were studied to find a catalyst with excellent activity and selectivity, which was not corrosive or destroyed by water.\nThere has been extensive work done with Ti catalysts, though most work occurred in conjunction with studies of Ziegler-Natta catalysts. The closest prior art known is U.S. Pat. No. 2,381,481 (Class 260-683.15), U.S. Pat. No. 2,951,885 (Class 260-671), U.S. Pat. No. 2,965,686 (Class 260-671) and U.S. Pat. No. 3,153,634 (Class 252-429).\nIn U.S. Pat. No. 2,381,481, preparation and use of a catalyst prepared by treating alumina gel with fluotitanic acid is disclosed. This catalyst is used for polymerization of olefins to heavier hydrocarbons, and also for alkylation of parafins with olefins, the latter when operating at high temperatures, between 700.degree. and 900.degree. F or higher. No mention is made of alkylation of aromatics with olefinic hydrocarbons or transalkylation of polyalkylbenzenes.\nIn U.S. Pat. No. 2,951,885, there is disclosed the use of titanium trihalide on activated alumina or other activated acidic oxide for alkylation of benzene with olefins. The catalyst is originally a tetrachloride, subsequently reduced to the trichloride with an alkali metal such as sodium, lithium, or potassium. The examples show that this catalyst will alkylate benzene with ethylene.\nIn U.S. Pat. No. 2,965,686 the thrust of the application was to develop a titanium subchloride catalyst. In Example II, a reaction between cumene and propylene was disclosed using titanium tetrachloride catalyst. The catalyst in Example II was tetrachloride. The catalyst was prepared by activating alumina by evacuation at a temperature of 600.degree. C for an unspecified period of time. The resultant catalyst was then used in an alkylation reaction for the propylation of cumene to form diisopropylbenzene. This patent is silent as to the type of alumina which was used as the base for the catalyst.\nIn U.S. Pat. No. 3,153,634, there is disclosed the use of titanium subhalides in a polymerization reaction. The patentee is probably describing a catalyst for production of solid polymer products. On page 3 lines 65-75, the patentee mentions use of benzene as an inert solvent to hold dissolved olefins, rather than as a reactant.\nAccordingly, work continued on developing an improved process for the catalytic conversion of aromatic hydrocarbons.\nAccordingly, the present invention provides a process for the catalytic conversion of an aromatic hydrocarbon comprising contacting the aromatic hydrocarbon with a reactant at aromatic hydrocarbon conversion conditions in the presence of a catalyst system comprising titanium tetrachloride and alumina wherein the catalyst system is prepared by passing TiCl.sub.4 vapor over activated alumina at a temperature of 20.degree. to 400.degree. C for 1 to 10 hours.\nThe catalyst used comprises titanium tetrachloride impregnated on an activated metal oxide which possesses surface hydroxyl groups. Specific examples of these metal oxides will include the Group III-A metal oxides which possess surface hydroxyl groups and which also possess a relatively high surface area such as alumina, gallium oxide, indium oxide, and thallium oxide. Of these compounds, the preferred substrate is alumina, and especially low density, high surface area aluminas such as gamma-alumina or, if so desired, eta-alumina.\nThe apparent bulk density of the alumina may range from about 0.3 to about 0.7 g/cm.sup.3 or higher with a surface area ranging from about 1 to about 500 m.sup.2 /g. The alumina may be in any shape, one example of the substrate being spheroidal alumina which is prepared by the conventional and commercial oil-drop method as described in U.S. Pat. No. 2,610,314. In addition, it is also contemplated within the scope of this invention that the alumina may be treated to provide greater physical stability, one type of treatment being to impregnate the gamma-alumina with a compound such as barium nitrate, which, upon calcination, is converted into barium oxide. The latter compound will then, as hereinbefore set forth, provide greater physical stability for the alumina. It is also contemplated within the scope of this invention that a commercial gamma-alumina may also be used as the support. However, since this commercial gamma-alumina could contain an excessive amount of water which would consume an excess of titanium tetrahalide without any beneficial effect on the catalyst, in the preferred embodiment of this invention the commercial gamma-alumina is subjected to a predrying step by heating to a temperature in the range of from about 400.degree. to about 550.degree. C under an inert gas or hydrogen flow for a period of about 1 to about 8 hours.\nIn order to achieve the maximum activity of the metal oxide support, it is necessary to avoid severe drying of said support. For example, the drying of alumina at temperatures in excess of about 600.degree. such as 650.degree., under vacuum, will seriously deplete the alumina of the hydroxyl groups present thereon. This severe drying step will not only remove the water which is absorbed on the alumina, but will also remove the aforesaid surface hydroxyl groups which are essential to make an active catalyst, said surface hydroxyl groups reacting with the titanium component of the titanium tetrachloride.\nThe gamma-alumina which has been predried according to the above paragraph is then placed in an appropriate apparatus which may comprise a flask, tube, etc., and a gas mixture of nitrogen and titanium tetrachloride which has been prepared by bubbling nitrogen gas through the liquid titanium tetrachloride at room temperature is passed over the gamma-alumina at temperatures of 25 to 135. Thereafter the temperature is increased to 550 or more. The passage of the nitrogen-titanium tetrachloride mixture over the alumina is effected about 0.5 to 10 hours or more, the time being dependent upon the amount of gamma-alumina which is present and the flow rate of the titanium tetrachloride-nitrogen gas mixture. It is preferred to pass the titanium tetrachloride the gaseous mixture over the support at a temperature of about 25.degree. to about 135.degree.. The temperature is then raised to 250 or a desired temperature, either gradually or in a series of steps. The preferred temperature for the heat treatment of this resulting composite is from about 135 to about 550, however, it will be dependent on the temperature which is used in the aromatic hydrocarbon conversion process. Generally, it is preferred that the temperature which is used in the treatment of the composite be equal to, or higher than, the aromatic hydrocarbon conversion temperature. Thereafter the temperature is maintained at this point and a stream of nitrogen is passed over the catalyst composite for an additional period which may range from about 1 to 10 hours. At the end of this time, the finished catalyst is then sealed under an inert atmosphere such as argon, helium, nitrogen, etc., prior to being used.\nAlternatively, the catalyst may be prepared by forming a solution of titanium tetrachloride in a polar, non-aqueous organic solvent and impregnating the alumina. Thereafter the impregnated alumina is then treated under a nitrogen flow at temperatures in the range hereinbefore set forth. After subjecting the impregnated substrate to these temperatures for a predetermined period of time, the finished catalyst is also recovered and maintained under an inert atmosphere until use thereof.\nIt is believed that the important factor in the thermal treating steps is the temperature, rather than the total time, as long as the total period for thermal treatment is reasonably long; around 5 or 6 hours. It is believed that the activity of the catalyst is effected by the thermal treatments because the thermal treatments desorb water molecules from the catalyst surface. Water can compete for active sites with the reactants, thus water is to some extent a catalyst poison. However, if catalyst deactivation occurs due to water adsorption, the catalyst can be regenerated by appropriate further thermal treatment under inert gas flow. Any regenerative thermal treatments should probably approximate those of the orignal thermal treatments, in that heating which is too rapid, or to too high a temperature, may cause hydrolysis of the TiCl.sub.4 component on the catalyst, which would reduce catalyst activity. Another danger of a rapid high temperature catalyst regeneration would be the formation of corrosive gases and liquids due to rapid evolution of H.sub.2 O vapor and chlorine compounds.\nIt is also contemplated within the scope of this invention that the catalyst system hereinbefore described may, if so desired, be composited on a solid support. The preferred solid supports which may be utilized comprise high surface area inert compounds, some representative examples of which will include silica, magnesia or mixtures of silica with other inorganic oxides such as silica-zirconia, silica-thoria, silica-magnesia-zirconia, etc., charcoal, coal, diatomaceous earths and clays such as fuller's earth, bentonite, montmorillonite, kieselguhr, etc. It is to be understood that these compounds will act only as supports for the catalyst system and will not enter into the catalytic activity of the composite. The catalyst system comprising titanium tetrachloride composited on the Group IIIA metal oxide and the inert support may be composited in any manner known in the art such as by impregnation, deposition, rolling, mixing, etc.\nIn addition, it is also to be considered within the scope of this invention that one or more promoters may be added to the catalyst system. It is believed that use of one or more promoters selected from the metals of Group VIB or Group VIII of the Periodic Table may be beneficial to the practice of the present invention.\nAt least about 0.5 weight percent titanium, on an elemental basis, is believed necessary for a significant amount of reaction to occur. The upper limit on titanium is believed to be about 20 wt. %.\nWhen it is desired to use the catalyst system in an alkylaromatic isomerization process, then alkylaromatic isomerization reaction conditions should be used. Reaction conditions are disclosed in U.S. Pat. No. 3,637,881 (Class 260-668a), the teachings of which are incorporated by reference. When it is desired to use the catalyst system of the present invention for alkylaromatic transalkylation then appropriate reaction conditions should also be used. These are disclosed in U.S. Pat. No. 3,720,726 (Class 260-672t), the teachings of which are incorporated by reference. Reaction conditions for the alkylation of aromatic hydrocarbons will be discussed in detail in a latter part of this specification.\nThe catalyst may be disposed in a reactor vessel as a fixed fluidized or moving bed of catalyst. The reactants may contact the catalyst in upflow, downflow or crossflow fashion, though upflow of reactants over a fixed bed of catalyst is preferred.\nThe liquid hourly space velocity in the reactor may range from 0.1 to 20. Because catalyst of the present invention is very active for the alkylation reaction, significantly higher space velocities are possible than when using some prior art catalysts, e.g., SPA. To some extent, the liquid hourly space velocity is related to temperature in the reaction zone, in general, a higher LHSV will require higher temperature operation.\nThe ratios of reactants and other reaction conditions which occur when alkylating benzene with light olefins, and preferably are basically those well known in the art. Pressures may range from 1 to 100 atmospheres, or even higher. It is desirable to maintain pressures high enough to have a liquid phase in the reaction zone. Although it is possible to operate at very high pressure, little advantage is gained thereby, in fact, an increase in pressure seems to have a harmful effect. Preferred pressure seems to be around 20 to 60 atm, with an optimum pressure of about 35 atm.\nTemperature effects both the conversion and selectivity of the reaction. Temperature may range between ambient and 250. At very low temperatures, the catalyst is not sufficiently active to permit the desired reaction to proceed at a satisfactory rate. At very high temperatures, it is believed that the catalyst may be damged, by formation of carbonaceous materials on the catalyst.\nIf the reaction is kinetically controlled, an increase in temperature should increase the rate of reaction. As a general statement, this is true, but the temperature dependence may not be as large as expected, if the reaction is limited by mass transport of reactants and products to and from the catalyst surface. Preferred operating temperature is probably about 100 to 200 C.\nThe catalyst may be disposed in a reactor vessel as a fixed, fluidized or moving bed of catalyst. The reactants may contact the catalyst in upflow, downflow or crossflow fashion, through upflow of reactants over a fixed bed of catalyst is preferred.\nThe liquid hourly space velocity in the reactor may range from 0.1 to 20. However, higher LHSV is possible depending on the desired conversion level of propylene. Because catalyst of the present invention is believed very active for the alkylation reaction, significantly higher space velocities should be possible than when using some prior art catalysts, e.g., SPA. To some extent, the liquid hourly space velocity is related to temperature in the reaction zone, in general a higher LHSV will require higher temperature operation."}
{"text": "1. Field of the Invention\nThe present invention relates to a high speed charge balancing comparator and, more particularly, to an improved charge balancing comparator having a substantially reduced response time whereby the speed of operation of the comparator is greatly improved.\n2. Description of the Prior Art\nCharge balancing comparators have been known and used for several years. The charge balancing comparator operates on the principle that the net charge onto its input nodes from two or more sources will cause either a positive or a negative voltage swing at the input charge balancing node. The state of the output of the comparator will depend on the polarity of the change of the voltage at this node. The comparator does not respond to the absolute value of the voltage at this node. At a suitable time, the output node is interrogated to see if it has swung in either a positive or a negative direction and this information is stored elsewhere.\nThis type of comparator is becoming increasingly popular in successive approximation analog-to-digital (A/D) converters using MOS or CMOS semiconductor integrated circuit technologies (although its use is not limited to these applications). One of the principle goals of any A/D converter is usually to minimize its time for making a conversion and most converters are ultimately limited in their high speed operation by the speed of their comparators.\nPrevious designs of comparators have been limited in speed of response due to the settling time of the input charge balancing node. The initial incorrect voltage on the input charge balancing node causes the output of the comparator, together with its internal nodes, to be disturbed significantly such that when the input charge balancing node does finally settle to its correct value, the comparator must first correct the internal node voltages, resulting in a long total response time at its output.\nA better understanding of this problem can be derived with reference to FIG. 1, which shows a schematic diagram of a typical charge balancing comparator, generally designated 10, and to FIG. 2, which shows a series of waveforms as a function of time which are useful in understanding the operation of comparator 10. For purposes of discussion, comparator 10 utilizes an operational amplifier 11 as the gain element. However, in practice, because of speed limitations, other configurations will almost always be used. Switches S.sub.1, . . . S.sub.N apply voltages V.sub.1, . . . V.sub.N, respectively, via capacitors C.sub.1, . . . C.sub.N, respectively, to input node A of operational amplifier 11. Similarly, switches S.sub.11, . . . S.sub.NN apply voltages V.sub.11, . . . V.sub.NN, respectively, via capacitors C.sub.1, . . . C.sub.N, respectively, to input node A. An autozero switch S.sub.AZ connects the inverting input of amplifier 11 to its output. Switches S.sub.1, . . . S.sub.N and S.sub.AZ are controlled by a timing signal shown as waveform 12 in FIG. 2 and switches S.sub.11, . . . S.sub.NN are controlled by a timing signal shown as waveform 13 in FIG. 2. Waveform 14 shown in FIG. 2 is the theoretical output V.sub.OUT of operational amplifier 11.\nDuring the autozero phase, switches S.sub.1, . . . S.sub.N and S.sub.AZ will be closed and switches S.sub.11, . . . S.sub.NN will be open. During this phase, switch S.sub.AZ connects the inverting input of amplifier 11 to its output, thereby putting amplifier 11 into a unity gain configuration. If amplifier 11 is perfect, the voltage V.sub.OUT will be at ground. In practice, because of finite errors such as input offset voltages, V.sub.OUT will be offset from ground potential up to a few millivolts. For the sake of simplicity, it will be assumed that amplifier 11 has no errors and that there are only two input pairs of voltages to be compared, V.sub.1 and V.sub.11 plus V.sub.2 and V.sub.22 (N=2 and NN=22). During the autozero phase, capacitors C.sub.1 and C.sub.2 will be charged to voltages V.sub.1 and V.sub.2, respectively.\nDuring the compare phase, switches S.sub.AZ, S.sub.1 and S.sub.2 are opened and switches S.sub.11 and S.sub.22 are closed. The voltage V.sub.A at the input node A to amplifier 11 will eventually assume a voltage V.sub.A : ##EQU1## It can be seen that V.sub.A can have any real positive or negative value depending on the magnitude of C.sub.1 and C.sub.2 and the magnitudes and polarities of (V.sub.11 -V.sub.1) and (V.sub.22 -V.sub.2).\nThe output voltage V.sub.OUT will respond after a finite time to the new value of V.sub.A. For example, if voltage V.sub.A has acquired a negative shift of value compared to that during the autozero phase, V.sub.OUT will swing by a positive amount as shown as waveform 14 in FIG. 2. At the end of this phase, V.sub.OUT is detected and stored elsewhere (on a capacitor or latch) and the cycle repeated. Each complete cycle is defined as the analog-to-digital conversion time and it is usually desired that this time be as short as possible.\nProblems arise in the use of comparator 10 when the mechanical switches shown in FIG. 1 are replaced by MOS switches or the like having finite values of series resistance and when the values in the numerator of equation (1) are different by a very small value to those in the denominator, i.e.: ##EQU2## The problem manifests itself in an initially inaccurate high value of voltage V.sub.A at node A at the beginning of the compare cycle. Although the voltage V.sub.A will finally settle to its correct value, this initial voltage can make the gain element (operational amplifier 11 in FIG. 1) swing its output and internal node voltages to incorrect values such that when the input voltage V.sub.A has settled down to near its correct value and final polarity, the gain element requires significantly more time to correct its output voltage V.sub.OUT and its internal node voltages compared to the case where this initial transient of input voltage had not occurred.\nTo demonstrate this initial transient, reference should be had to FIG. 3, wherein the input of comparator 10 is modified by eliminating switches S.sub.11 and S.sub.22 and substituting therefore single switches S.sub.1 and S.sub.2 with their respective series resistors R.sub.1 and R.sub.2. Additionally, for the sake of simplicity, switch S.sub.AZ is shown connecting node A to ground during the autozero phase. Switches S.sub.1 and S.sub.2 also connect capacitors C.sub.1 and C.sub.2, respectively, to voltages V.sub.1 and V.sub.2, respectively, during the autozero phase.\nThe voltage at node A from the beginning of the compare phase is defined by the equation: EQU V.sub.A =V.sub.A(FINAL) +[V.sub.A(INITIAL) -V.sub.A(FINAL) ]e.sup.-A/t, (2)\nwhere t=time, ##EQU3##\nIt should be noted that the initial value of V.sub.A is solely dependent upon the input differential voltages (V.sub.11 -V.sub.1) and (V.sub.22 -V.sub.2) and the equivalent resistance of R.sub.1 and R.sub.2, that the final value of V.sub.A is solely dependent on the same differential voltages and capacitors C.sub.1 and C.sub.2, and that the time constant for the voltage V.sub.A to change from V.sub.A(INITIAL) towards V.sub.A(FINAL) is dependent on the equivalent RC product of the two resistors R.sub.1 and R.sub.2 in series and the two capacitors C.sub.1 and C.sub.2 in series.\nTo demonstrate this situation, consider a practical case for an 8-bit A/D converter. The reference voltage (V.sub.11 -V.sub.1) is -2.500 volts and the full scale input voltage (V.sub.22 -V.sub.2) less one least significant bit is (5.000-5.000/256) volts or approximately 4.980 volts. In this example, capacitor C.sub.1 has a value of 10 picofarads, capacitor C.sub.2 has a value of 5 picofarads and the equivalent series resistances R.sub.1 and R.sub.2 of the MOS switches are 2000 ohms each. Inserting these values in equations (5) and (1): ##EQU4## The time constant TC for V.sub.A is defined by: ##EQU5## Between four and five time constants (approximately 14 nanoseconds) are required before voltage V.sub.A assumes approximately its correct value.\nThis example is shown in FIG. 4 where voltage V.sub.A is plotted as a function of time (curve 15). The ratio of the initial magnitude of voltage V.sub.A is almost 200 times that of its final value. More importantly, its initial polarity is reversed from its final condition. As the example shows, the initial input transient can be of such a high magnitude and sign that it can cause the gain element of the comparator to slam its output into an incorrect state and, often just as undesirable, cause its internal nodes to assume grossly incorrect voltage values such that when the input voltage V.sub.A assumes its correct polarity, at time t.sub.1 in FIG. 4, and then its small terminal magnitude, the internal node and output voltages have to be initially corrected to their autozero values and then to their final values. This produces a delay at the output of the comparator which is usually twice to many times that which would occur if no input transient had been allowed and the input voltage V.sub.A had been connected to the input of the comparator only after the time t.sub.1. Since the speed or response time of the comparator is the limiting factor for A/D converters and other systems using comparators, it is highly desirable to eliminate or reduce the effect of input voltage transients for this family of charge balancing comparators."}
{"text": "(1) Field of the Invention\nThis invention relates to a technology for making smaller and lighter RF passive circuits and RF amplifiers equipped with via-holes.\n(2) Prior Art\nRecently, various types of mobile communication tools, such as portable phones or portable information terminals have been commercialized all over the world. As portable phones, cellular phones for bands of 900 MHz and 1.5 GHz, and Personal Handyphone System (PHS) for a band of 1.9 GHz are two examples that are commercialized in Japanese market. Other examples include world-famous GSM, and CDMA among the technologies adopted in PCS (Personal Communications Services) in the U.S.A.\nAs a third-generation mode following the analogue mode and the digital mode, IMT2000 is planned to be commercialized in the future.\nIn developing mobile communication terminals especially portable terminals, it is an inevitable trend to seek smaller and lighter terminals. Accordingly, it is important to achieve a technology for making smaller and lighter components for these terminals.\nAs a trend, it is desired to make high frequency components of the portable terminals as a monolithic microwave IC (MMIC). The MMIC, in which active elements, their matching circuits, and bias circuits are integrated on the same substrate, is more advantageous in making smaller products than a Hybrid IC which is structured to have circuits and bias electricity-feeding circuits as outside-chips.\nEven using the MMIC, it is required to ground circuit elements. Conventional grounding methods include a method of wire-bonding from the surface of semiconductor substrates, and a via-hole method. It is more effective to use the via-hole method in achieving high-quality and low cost for packaging, which makes the via-hole method more frequently adopted in the MMIC.\nThe following is a description of an example of a conventional type of RF passive circuit and RF amplifier equipped with via-holes with reference to FIGS. 8A–8D.\nFIG. 8A is a schematic circuit diagram of a conventional RF amplifier which includes RF passive circuits equipped with via-holes, and FIGS. 8B and 8C are pattern diagrams of conventional RF passive circuits both equipped with a via-hole.\nAs FIG. 8A shows, a source-ground type of RF amplifier is constructed by connecting: a gate bias resistance 805 and an input matching circuit 806 to a gate terminal 802; a drain voltage feeding circuit 807 and an output matching circuit 808 to a drain terminal 803; and a source terminal 804 to a ground terminal 809, in the field effective transistor (FET) 801. An input terminal 810 and an output terminal 811 are both 50 Ω impedance, and the input matching circuit 806 and the output matching circuit 808 are adjusted to 50 Ω. Further, each of an input DC cut capacitor 812 and an output DC cut capacitor 813 is inserted to the input side and the output side respectively.\nThe input matching circuit 806 consists of an input matching parallel inductor 814, an input matching parallel capacitor 815, and an input matching serial inductor 816. The input matching parallel capacitor 815 is grounded by an input matching circuit via-hole 821.\nThe output matching circuit 808 consists of an output matching serial inductor 817, and an output matching parallel capacitor 818. The output matching parallel capacitor 818 is grounded by an output matching circuit via-hole 822.\nThe drain voltage feeding circuit 807 consists of a choke inductor 819 and a bypass capacitor 820. The bypass capacitor 820 is grounded by a drain voltage feeding circuit via-hole 823.\nFIGS. 8B and 8C are both pattern diagrams of an RF passive circuit with a via-hole; each of them shows the input matching circuit 806 and the drain voltage feeding circuit 807 respectively. FIG. 8D shows a cross-sectional view taken along line (A–A′) of FIG. 8B. The following is a description of a common part between the input matching circuit 806 and the drain voltage feeding circuit 807, taking an example of the input matching circuit 806.\nConstituting elements of the aforementioned input matching circuit 806 is made, as a semiconductor substrate, on a surface of a GaAs substrate 824. Both of the input matching parallel inductor 814 and the input matching serial inductor 816 are made in a spiral-electrode-pattern, and the input matching parallel capacitor 815 is made in an MIM (Metal-Insulator-Metal) capacitor pattern.\nAs FIG. 8D shows, the spiral-electrode-pattern is made on the GaAs substrate 824 which is covered by an insulator film 834 such as silicon oxide. Specifically, the spiral-electrode-pattern is a structure where a lower wiring metal layer 831 which is made by gold/titanium vacuum evaporation is connected to an upper wiring metal layer 830 made by gold-plating by means of a contact hole 833, with a between-layer insulator film 832 in between.\nOn the other hand, the MIM capacitor is a structure where an upper wiring metal 829 is formed on a dielectric layer 828 under which is an electrode extended from the lower wiring metal layer 831; the upper wiring metal 829 is made by gold/titanium vacuum evaporation and the dielectric layer 828 is titanium oxide strontium (SrTiO3:STO) with a permittivity of 100 or more. The end of the electrode extended from the upper wiring metal 829 is connected to a ground metal layer 826 which is situated on the via-hole, as FIGS. 8B and 8C show.\nThe input matching circuit via-hole 821 can be formed by etching from the main surface of the GaAs substrate 824 where circuit elements were made (a surface via-hole). Or, it could also be formed by etching from the other main surface (a backside via-hole). Inside the via-hole 821, an electric conducting film is conducted to a backside ground metal 829. This electric conducting film is electrically connected to the upper wiring metal 829 of the MIM capacitor through the ground metal layer 826.\nFurther, as depicted in FIG. 8C, constituting elements of the drain voltage feeding circuit 807 are formed, as a semiconductor substrate, on the surface of the GaAs substrate 824. As for the choke inductor 819, a spiral-electrode-pattern is used, and as for the drain voltage feeding circuit via-hole 823, either a surface via-hole or a backside via-hole is used for forming.\nNote that a feeding terminal 825 is structured by extending a drain voltage terminal 836 from the lower wiring metal layer 831 through an extending wire 835.\nThus structured as above, the following constituting elements of the RF passive circuit are formed on and through the GaAs substrate: the spiral inductor, the MIM capacitor, and the via-hole. Moreover, as FIG. 8C shows, the above three elements are positioned at a different location two-dimentionally, and are connected to each other by wiring. The elements constitute the RF amplifier with a help of the input matching circuit 808 and the drain voltage feeding circuit 807.\nAs seen above, the conventional type of RF amplifiers and RF passive circuits cannot be made smaller in size, due to the two-dimensional positioning of the constituting elements of the drain voltage feeding circuit 807, which inherently take much space."}
{"text": "For various apparel articles such as jackets, baseball caps, and other products, as well as accessory articles such as backpacks and cinch bags, it may be desirable to have adornments, such as a patch, attached thereto. In certain applications use of a removable adornment (e.g. a removable patch) is advantageous as the indicia on the apparel or accessory can be changed by the user by simply replacing the adornment.\nA removable adornment can be attached to apparel and other products and accessories (articles) via numerous fastening techniques. For example, a patch can be removably attached to an article of manufacture using Velcro® (a hook-and-loop fastener or attachment system), magnets, adhesives, snaps, buttons, and pins. However, such attachment techniques often result in the patch protruding away from the surface of the apparel or product, which can look low-cost, unprofessional, and/or permit the patch to be unintentionally or accidently removed. In addition, the entire surface of the patch including its borders may not be flush with the surface of the apparel or accessory thereby causing wrinkles or other unevenness of the patch and its borders, again resulting in a low-cost, unsophisticated look.\nIn certain designs, the patch can have a border such as a stitched border around the patch. The stitched border can be a merrowed border or a satin stitched border. Although such borders on a patch can improve the look of the patch on the apparel or accessory, similar issues can be present for removable patches with borders as with removable patches without borders (e.g., patches having hot cut borders).\nThus, there is a need to improve the removable attachment of adornments to apparel and other products and accessories, which attachment can include patch receptacles for removable patches, such that imperfections of the adornment cannot be readily seen and/or the resulting combination of the adornment and the article can present the look of a bordered adornment on the apparel or accessory, despite the adornment being borderless."}
{"text": "The invention is based on a fuel injection pump for internal combustion engines having a cam drive effecting the supply movement of at least one pump piston, the cam drive having a relatively rotatable portion supported in a pump housing for adjustment of injection onset, an adjuster piston cooperating with the cam drive and subjected to the rpm-dependent pressure of a supply pump counter to the force of two restoring springs. In known fuel injection pumps of this kind, either both restoring springs function simultaneously over the entire rpm range, or the action of one spring is added to that of the other once the adjusting piston has traveled a certain portion of its stroke. In the first case, the selection of different characteristic curves for the springs results in an overall characteristic curve which flattens out in the upper rpm range; however, there are limits to this flattening. In the second case, postponing the functioning of the second spring results in a clearly defined break in the characteristic curve, although with the disadvantage that the only control sequences which are possible are those where initially, at low rpm (that is, until the second spring comes into engagement), relatively large relative rotations per change in the rpm are possible and subsequently only smaller relative rotations per rpm chamber. However, in many internal combustion engines it is necessary for the course of the characteristic curve for injection onset to be reversed; that is, the injection onset, beyond a specific rpm, should adjust more rapidly toward \"early\" per rpm change with increasing rpm than is desired in the same engine at lower rpm. The instant of injection onset has an increasingly important role, given increasingly stringent demands for smoothness of engine operation and for the nontoxicity of exhaust gases."}
{"text": "The process of blood clotting begins with an injury to a blood vessel. The damaged vessel wall initiates hemostasis by causing adherence and accumulation of platelets at the point of vascular injury and by activating the plasma proteins which initiate the coagulation process. Sequential activation, via specific proteolytic cleavages and conformational changes, of a series of proteins (including Factor VIII) initiate a coagulation cascade that eventually leads to deposition of insoluble fibrin. The fibrin, together with aggregated platelets, curtails the escape of blood through the point of injury in the damaged vessel wall.\nOne of these coagulant proteins, Factor VIII, is a plasma protein that has the ability to correct the clotting defect in plasma from patients with Hemophilia A. The activity of Factor VIII is also measured by its ability to induce clotting in plasma obtained from these patients. One unit of Factor VIII is defined as the amount of Factor VIII present in one milliliter of normal, adult male plasma. This standard is a World Health Organization standard and is available from the National Institute for Biological Standards and Control, Holly Hill, Hampstead, London NY3 6RD, England.\nDuring normal blood coagulation, Factor VIII is activated by thrombin, a protease, which affects a proteolytic modification (cleavage) of Factor VIII. (Weiss et al., Science, 182: 1149-1151, 1973; Hoyer, Chapter 4, Hemostasis and Thrombosis, 2d Ed., Colman, Ed., J. B. Lippincott & Co., Philadelphia, Pa., 1987; Fulcher et al., J. Clin. Invest., 6: 117-124, 1985; and Eaton et al., Biochemistry, 25: 505-512, 1986). Following an initial increase in activity after exposure of Factor VIII to thrombin, there is a decay in activity to below base line levels. This decay in activity is attributed to further proteolytic cleavage of activated Factor VIII by thrombin, instability of the partially cleaved molecules Hultin, M. B. et al., Blood, 57: 476-482, 1981, and during pathological states, to inactivation by thrombin-activated protein C (Eaton, D. et al., Bioch. 25: 505, 1986). The activation and inactivation of Factor VIII by thrombin and protein C allows a fine measure of control over coagulant activity in vivo. Unfortunately, it may also cause premature degradation and instability of Factor VIII, and particularly of the isolated, purified Factor VIII or purified plasma fractions containing Factor VIII, prior to its administration to mammalian hosts in need of such treatment.\nThe primary therapeutic use of Factor VIII has been its intravenous administration to patients with hemophilia A to stop or to prevent bleeding. Early administration was by the infusion of whole blood or fresh-frozen plasma. These methods required the infusion of very large volumes of proteins and long infusion times, often causing hypervolemia.\nThe use of plasma cryoprecipitate (\"cryo\") eliminated the need for whole blood or fresh frozen plasma infusion but the moderate volumes and high protein concentrations still required long infusion times. The cryo was not completely soluble, the cryo solution used for infusion required filtration, and the cryo's AHF content was low and highly variable. Moreover, cryo required storage at -20.degree. C. in large plastic blood bags, and therefore, required large refrigeration facilities.\nThe primary source of native Factor VIII used today to treat those in need of clotting factor is still plasma, and the use of lyophilized, higher purity Factor VIII concentrates has solved many of the above problems. However, the yield obtained by conventional methods of purification is relatively low, making the cost of the factor excessive. Furthermore, some of the low purity concentrates contain high amounts of contaminating proteins including specific blood-type antibodies that cause hemolysis, proteins that cause immunological abnormalities including a temporary inversion of the T-cell ratio (helper/suppressor) which resembles AIDS (Carr, R. et al., Lancet 1: 1431, 1984), fibrinogen, fibronectin, Von Willebrand factor, and other proteins.\nWhen a clinical Factor VIII concentrate shows evidence of prior activation, several in vitro and in vivo problems may occur. If the Factor VIII assay was performed by a one stage procedure, it may be falsely high (Niemitz and Nossel, Brit. J. of Hematol. 16: 337, 1969); this is not the case, if the assay is performed by a two stage procedure, but very few laboratories are equipped to use the method. Additionally, the activated material is unstable and tends to lose activity during short term storage in the liquid state and during long term storage in the lyophilized state, resulting in poor results when infused in vivo. Use of these Factor VIII concentrates could lead to inadequate treatment for the patient and excessive bleeding.\nSubstantial activation of antihemophilic factor in plasma and cryoprecipitate is usually prevented by the addition of citrate as the blood is drawn. However, the addition of too much citrate at the time of drawing or later in the purification procedure can cause inactivation of Factor VIII and a corresponding decrease in yield (Rock, et al., Thromb. Res. 29: 521-535, 1983, Foster, P. R. et al., Scan. J. Haematol. Supp. 40, 33: 103-110, 1984). Consequently, heparin and calcium are often added to the cryoprecipitate or the partially purified Factor VIII during fractionation. In addition, most of the contaminating procoagulant clotting factors or proproteases which are present, such as Factor II, VII, IX and X, and protein C, are adsorbed with aluminum hydroxide. However, it is difficult, if not impossible, to adsorb all of the contaminating Factors II, VII, IX, and X and protein C with aluminum hydroxide without also adsorbing Factor VIII. For example, when Factors II, VII, IX and X are adsorbed from cryoprecipitate or plasma with 5 to 20% aluminum hydroxide and the residual cryoprecipitate or plasma is brought to approximately 3 mM Ca.sup.+ , coagulation still occurs in a relatively short period. It was believed that this was due to residual Factors II, VII, IX, and X in the preparation.\nAdditionally, minute amounts of thrombin tend to form during blood collection and during large scale cryoprecipitation. When Factor VIII is purified and concentrated, some thrombin is purified and concentrated in addition. The effects of various proteolytic enzymes, including thrombin on purified Factor VIII, were shown by Eaton et al. Bioch. 25: 505 (1986), and the activation of Factor VIII by thrombin with redistribution of polypeptide chains in Factor VIII was demonstrated by polyacrylamide gel electrophoresis by Fulcher et. al., Blood 61: 807-811, (1983), Vehar et al., Nature, 312, 337-342 (1984), and Rotblatt et al., Biochem. 24: 4294-4300 (1985). Similar biochemical changes were found in commercial Factor VIII concentrates purified by large scale methods for clinical use when analyzed by polyacrylamide gradient autoradiograms of Factor VIII-Factor VIII antibody complexes (Barrowcliffe, Lancet, 1: 1448-1449 (1986).\nChromatographic techniques, including ionic, hydrophobic, and immunoaffinity chromatography, have been used to purify Factor VIII, but the yield is relatively low. Several groups of investigators have purified Factor VIII with specific activities as high as 3,000 units per milligram of protein (Fulcher et al., Proc. Natl. Acad. Sci. USA, 79: 1648-1652, (1982); Tuddenham, et al., NIH Symposium on Factor VIII/von Willebrand Factor, pg. 1., Scripps Clinic and Research Foundation, La Jolla, Calif. (1982); Johnson et al., NIH Symposium on Factor VIII/von Willebrand Factor, pg. 2., Scripps Clinic and Research Foundation, La Jolla, Calif. (1982).\nMuch effort has been made to increase the purity of Factor VIII concentrates. This included the use of immunoaffinity chromatography and monoclonal antibodies (Fulcher et al., Proc. Natl. Acad. Sci. USA. 79: 1648-1652 (1982) and Zimmerman et al., U.S. Pat. No. 4,361,509. Overall recovery was estimated to be less than 15% without a heating step.\nFurther attempts were described by J. J. Morgenthaler, Thromb. Haemostas, 47(2): 124 (1982) wherein the use of acetate-glycine-lysine buffer on modified sepharose columns to purify Factor VIII:C from polyethylene glycol precipitated AHF was disclosed.\nThe addition of 0.05M dextrose in the eluting buffer of bovine Factor VIII during ion exchange column chromatography (on DEAE-cellulose) slightly improved the purity of the product and increased the yield as reported by Lundblad et al. Thrombosis Research, 1: 197, Pergamon Press, Inc. (1972). However, resolution was sightly reduced.\nPurification of Factor VIII by column chromatography in the presence of additives, including sugars and polyhydric alcohols which serve to increase the electrostatic forces on the surface of the proteins while decreasing the hydrophobicity of the proteins, was reported to result in high recovery of preparations having high purity and resolution (Mathews et al., U.S. Pat. Nos. 4,743,680, 4,847,362, and 4,952,675).\nFulton et al., U.S. Pat. No. 4,970,300, disclose conjugates of proteins having antihemophilic factor activity linked to non-antigenic ligands. These conjugates were stated to have a longer half-life than unconjugated proteins while maintaining substantial Factor VIII activity. Fulton et al. also used short columns of ionic resins to remove residual amounts of Factors II, VII, IX and X.\nFinally recombinant Factor VIII has become available through the cloning of the Factor VIII gene. (Gitschier, et al. Nature 312: 326-330 (1984); Wood et al., Nature 312: 330-337, 1984; Vehar et al., Nature, 312: 337-342, (1984); and Toole, et al., Nature, 312: 342, (1984).\nA novel method for stabilizing partially purified Factor VIII and Factor VIII from biological fluids or mixtures of cellular products derived from recombinant materials, often containing calf serum and/or endopeptidases, has now been discovered. This method prevents activation or degradation of moderately, or highly purified Factor VIII during viral inactivation, purification, recovery, lyophilization and storage and produces a higher yield from the starting material in vitro and a higher recovery in vivo. This results, in turn, in a safer, more effective therapeutic agent at lower cost."}
{"text": "Silicon-based integrated circuits have long been used as a platform for microelectronic applications. For example, microprocessors in computers, automobiles, avionics, mobile devices, control and display systems and in all manner of consumer and industrial electronics products are all traditionally based on a silicon platform that facilitates and directs the flow of electricity. As processing requirements have increased, the design of silicon-based integrated circuits has adapted to accommodate for faster processing times and increased communication bandwidths. Primarily, such performance gains have been the result of improvements in feature density, meaning that technologies have been developed to crowd ever-increasing numbers of features such as transistors onto a silicon chip. While efforts to increase feature density continue, alternative methods for increasing processing speeds and bandwidth on silicon-based platforms are also being developed. One such method is known as silicon photonics.\nThe term “silicon photonics” relates to the study and application of photonic systems that use silicon as an optical medium. Thus, instead of or in addition to using silicon to facilitate the flow of electricity, silicon is used to direct the flow of photons or light. While the speed of electricity and the speed of light are the same, light is able to carry data over a wider range of frequencies than electricity, meaning that the bandwidth of light is greater than that of electricity. Thus, a stream of light can carry more data than a comparable stream of electricity can during the same period of time. Accordingly, there are significant advantages to using light as a data carrier. Furthermore, using silicon as a preferred optical medium allows for application of and tight integration with existing silicon integrated circuit technologies. Silicon is transparent to infrared light with wavelengths above about 1.1 micrometers. Silicon also has a high refractive index of about 3.5. The tight optical confinement provided by this high index allows for microscopic optical waveguides, which may have cross-sectional dimensions of only a few hundred nanometers, thus facilitating integration with current nanoscale semiconductor technologies. Thus, silicon photonic devices can be made using existing semiconductor fabrication techniques, and because silicon is already used as the substrate for most integrated circuits, it is possible to create hybrid devices in which the optical and electronic components are integrated onto a single microchip.\nIn practice, silicon photonics are implemented using silicon-on-insulator, or SOI, technology. In order for the silicon photonic components to remain optically independent from the bulk silicon of the wafer on which they are fabricated, it is necessary to have an intervening material. This is usually silica, which has a much lower refractive index of about 1.44 in the wavelength region of interest. This results in total internal reflection of light at the silicon-silica interface and thus transmitted light remains in the silicon.\nA typical example of data propagation using light is illustrated in FIG. 1. FIG. 1 illustrates an optical transmission system 100 that includes, for example, a silicon waveguide 110. The silicon waveguide may make up the entirety of the optical transmission system 100 or just one or more portions of the system 100. The system includes multiple data input channels 120, where each channel 120 transmits data in the form of pulses of light. In order to simultaneously transmit the data carried on the multiple data channels 120, the light in each channel 120 is modulated by a frequency modulator 130. The modulated light from each channel 120 is then combined into a single transmission channel 150 using an optical multiplexer 140. The multiplexed light is then transmitted along the single transmission channel 150 to an endpoint (not shown) where the light is de-multiplexed and demodulated before being used by an endpoint device.\nTransmission of light in an optical waveguide is, however, affected by temperature. In general, changes in temperature can result in changes in the device dimensions (due to thermal expansion) and refractive indices of the materials used in the optical waveguide. More particularly, changes in temperature can affect the operation of the optical frequency modulators 130 illustrated in FIG. 1. Resonant photonic modulators are designed to only modulate received frequencies that are at or close to specific known frequencies. To only allow the modulation of the specific known frequencies, the modulators include resonant structures that act to filter out all but the known frequencies which are to be modulated by the modulators. Thus, the known frequencies are resonant frequencies of the resonant structures. Unfortunately, because the refractive indices of the resonant structures tend to change according to temperature, the specific frequencies that are modulated (i.e., the resonant frequencies) tend to deviate from the known frequencies as the temperature changes. Therefore, there is a need for silicon optical waveguides with modulator circuits that are tolerant of changes in temperature."}
{"text": "The invention relates to weight transfer devices for tractor trailer vehicles, and more particularly to a weight transfer device placed behind a tractor's fifth wheel for exerting an upward pressure against a trailer underside for forcing down the rear tandem drive wheels of the tractor so as to maintain or enhance their traction for example on uneven ground where the trailer tends to uplift the drive wheels of the tractor."}
{"text": "This invention relates to electronic systems. In particular, the present invention is a backplane assembly for an electronic system chassis, wherein the backplane assembly incorporates connectors for a variety of printed circuit board modules in a multiple channel architecture.\nThe general construction and operation of electronic systems for an electronic device necessitates that a plurality of printed circuit board modules, such as input/output (I/O) modules be electronically connected to one or more central processing units. In order to achieve this electronic connection, each I/O module electronically contacts a mother board, often referred to as a backplane board as it is vertically mounted as a back wall of an electronic system chassis. The backplane includes a plurality of connectors for releasably interconnecting with a mating connector along a rear edge of each of the I/O modules. Typically, a pair of opposed side walls of the chassis incorporate guide channels for providing lateral support to the I/O modules when interconnected with the backplane.\nFor fail/safe operation, as may be necessary in electronic devices of aircraft systems, construction of the electronic system necessitates the use of redundant I/O modules. Typically, to achieve such redundancy requires the use of two chassis, with each chassis having an identical set of I/O modules. Each backplane of the redundant chassis would be linked to the electronic device through its own signal transmitting bus or channel. In other words, I/O modules of one chassis would be linked to the electronic device through a so called \"channel A\", while the redundant, second chassis would be linked to the electronic device through a so called \"channel B\". Operational failure of an I/O module within the one chassis would only require switching from the inoperative channel A to the operative channel B of the operational redundant, second chassis, to maintain the operation of the electronic device. One drawback of this redundant chassis architecture, particularly with regard to aircraft where space is at a premium, is the space requirement for the chassis of the redundant, second set of I/O modules.\nU.S. Pat. No. 5,547,386 to Fredberg in one way, addresses this redundant chassis, space requirement drawback of redundant channel electronic systems. In Fredberg, a single chassis, having a single backplane incorporating redundant sets of connectors is used. The backplane connectors for the I/O modules are arranged side by side, with adjacent backplane connectors alternating between channel A and a channel B. Though the backplane architecture of Fredberg eliminates the added bulk of the redundant, second chassis, there are some disadvantages to the Fredberg design. For example, since only a single backplane is used, it is difficult to discern which I/O modules are on channel A and which are on channel B. This difficulty is particularly evident when inoperative or damaged I/O modules of an inoperative channel need to be identified for repair or replacement. In addition, due to the alternating nature of the channel A and channel B backplane connectors of Fredberg, I/O modules of a common channel can not be arranged side by side. Moreover, the backplane scheme of Fredberg can not accommodate \"ganging\" (i.e., connecting multiple I/O modules directly together) of like I/O modules.\nThere is a need for improved redundant channel, backplane assemblies for chassis of electronic systems. In particular, there is a need for a redundant channel, backplane assembly wherein I/O modules on channel A can be quickly discerned from I/O modules on channel B. This would allow I/O modules of an inoperative channel (inoperative due to inoperative or damaged I/O modules) to be readily identified with respect to the operative I/O modules of the operative redundant channel. In addition, the redundant channel, backplane assembly should permit I/O modules of a common channel to be arranged side by side, and should accommodate ganging of like I/O modules. Lastly, the redundant channel, backplane assembly should be as compact as possible so as to be particularly suitable for use in aircraft or wherever space is at a premium."}
{"text": "1. Field of the Invention\nThe present invention relates to cleaning compounds for cleaning insect splatters from a windshield and more particularly pertains to a new cleaning compound additive for adding to water or windshield wiper fluid to aid removal of insect splatters from a windshield.\n2. Description of the Prior Art\nThe use of cleaning compounds for cleaning insect splatters from a windshield is known in the prior art. More specifically, cleaning compounds for cleaning insect splatters from a windshield heretofore devised and utilized are known to consist basically of familiar, expected and obvious configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nKnown prior art includes U.S. Pat. No. 4,507,155; U.S. Pat. No. 5,711,956; U.S. Pat. No. 3,607,759; U.S. Pat. No. 3,337,466; U.S. Pat. No. 2,931,776; and U.S. Pat. No. Des. 301,637.\nWhile these compositions fulfill their respective, particular objectives, and requirements, the aforementioned patents do not disclose a new cleaning compound additive. The inventive compound includes the combination of effective amounts of anhydrous ammonia, alkoxylated linear alcohols sold under the trade name SURFONIC, dry particulate sodium bicarbonate, dry particulate polyvinyl pyrrolidone, dry particulate orange dye, and dry particulate citrus orange scent.\nIn these respects, the cleaning compound additive according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in so doing provides an compound primarily developed for the purpose of adding to water or windshield wiper fluid to aid removal of insect splatters from a windshield.\nIn view of the foregoing disadvantages inherent in the known types of cleaning compounds for cleaning insect splatters from a windshield now present in the prior art, the present invention provides a new cleaning compound additive construction wherein the same can be utilized for adding to water or windshield wiper fluid to aid removal of insect splatters from a windshield.\nThe general purpose of the present invention, which will be described subsequently in greater detail, is to provide a new cleaning compound additive composition and method which has many of the advantages of the cleaning compounds for cleaning insect splatters from a windshield mentioned heretofore and many novel features that result in a new cleaning compound additive which is not anticipated, rendered obvious, suggested, or even implied by any of the prior art cleaning compounds for cleaning insect splatters from a windshield, either alone or in any combination thereof.\nTo attain this, the present invention generally comprises the combination of effective amounts of anhydrous ammonia, alkoxylated linear alcohols sold tinder the trade name SURFONIC, dry particulate sodium bicarbonate, dry particulate polyvinyl pyrrolidone, dry particulate orange dye, and dry particulate citrus orange scent.\nThere has thus been outlined, rather broadly, the more important features of the invention in order that the detailed description thereof that follows may be better understood, and in order that the present contribution to the art may be better appreciated. There are additional features of the invention that will be described hereinafter and which will form the subject matter of the claims appended hereto.\nIt is to be understood that the phraseology and terminology employed herein are for the purpose of description and should not be regarded as limiting.\nAs such, those skilled in the art will appreciate that the conception, upon which this disclosure is based, may readily be utilized as a basis for the designing of other compositions, methods and systems for carrying out the several purposes of the present invention. It is important, therefore, that the claims be regarded as including such equivalent compositions insofar as they do not depart from the spirit and scope of the present invention.\nFurther, the purpose of the foregoing abstract is to enable the U.S. Patent and Trademark Office and the public generally, and especially the scientists, engineers and practitioners in the art who are not familiar with patent or legal terms or phraseology, to determine quickly from a cursory inspection the nature and essence of the technical disclosure of the application. The abstract is neither intended to define the invention of the application, which is measured by the claims, nor is it intended to be limiting as to the scope of the invention in any way.\nIt is therefore an object of the present invention to provide a new cleaning compound additive compounds and method which has many of the advantages of the cleaning compounds for cleaning insect splatters from a windshield mentioned heretofore and many novel features that result in a new cleaning compound additive which is not anticipated, rendered obvious, suggested, or even implied by any of the prior art cleaning compounds for cleaning insect splatters from a windshield, either alone or in any combination thereof.\nIt is another object of the present invention to provide a new cleaning compound additive which may be easily and efficiently manufactured and marketed.\nAn even further object of the present invention is to provide a new cleaning compound additive which is susceptible of a low cost of manufacture with regard to both materials and labor, and which accordingly is then susceptible of low prices of sale to the consuming public, thereby making such cleaning compound additive economically available to the buying public.\nStill another object of the present invention is to provide a new cleaning compound additive for adding to water or windshield wiper fluid to aid removal of insect splatters from a windshield.\nYet another object of the present invention is to provide a new cleaning compound additive which includes the combination of effective amounts of anhydrous ammonia, alkoxylated linear alcohols sold under the trade name SURFONIC, dry particulate sodium bicarbonate, dry particulate polyvinyl pyrrolidone, dry particulate orange dye, and dry particulate citrus orange scent.\nStill yet another object of the present invention is to provide a new cleaning compound additive that enhances the ability for windshield wiper fluid to remove insect spatters from a windshield."}
{"text": "There are many conventional devices used in the field of children's books to attract and maintain an attention span of children. A commonly used conventional device provides a large, colorful illustrations showing the subject matter (e.g., a story) of the book or providing an eye-catching and amusing background for the story.\nAnother such conventional device is described in U.S. Pat. No. 4,819,963, which describes a book sculptured in the shape of a fish or other form, apparently for the purpose of gaining the interest of children. However, it would be desirable to further attract the interest of children in order to encourage children to read."}
{"text": "Retroreflective appliques have the ability to return a substantial portion of incident light back towards the light source. This unique ability has promoted widespread use of retroreflective appliques on clothing. Persons who work or exercise near motor vehicle traffic need to be conspicuously visible so that they do not get struck by passing motor vehicles. When placed on clothing, the retroreflective appliques highlight a person's presence by retroreflecting light from motor vehicle headlamps.\nA retroreflective applique typically comprises an optical lens element layer, a polymeric binder layer, and a reflective layer. The optical lens elements commonly are microspheres that are partially embedded in the polymeric binder layer. The reflective layer typically comprises aluminum, silver, or a dielectric mirror that usually is disposed on the embedded portions of the microspheres. Light striking the front surface of the retroreflective applique passes through the microspheres and is reflected by the reflective layer to re-enter the microspheres where the light's direction is then altered to travel back towards the light source.\nRetroreflective appliques have been made by partially embedding a microsphere layer in a thermoplastic carrier web, applying a reflective material over the microspheres' protruding portions, and then forming a binder layer over the coated microspheres. Often a pressure-sensitive adhesive is applied on the binder layer's back surface, and a release liner is placed over the adhesive until the applique is secured to a substrate. The completed applique (also referred to as a transfer sheet) is supplied to a garment assembler in this form, and the garment assembler secures the applique to an article of clothing by removing the release liner and adhering the applique to an outer surface of the article of clothing. The carrier is then separated from the applique to expose the microspheres so that the applique can retroreflect light. Although known retroreflective appliques demonstrate very good retroreflective performance and are very effective in highlighting a wearer's presence when light strikes the clothing at nighttime, the known appliques possess some drawbacks.\nA first drawback is that the applique requires use of a carrier web that becomes discarded as waste by the garment assembler. The carrier web is not separated from the applique before the manufacturer supplies the applique to the garment assembler because the applique is not yet secured to a substrate. If the applique was separated from the carrier beforehand, the applique's binder layer and the reflective layer can become irreversibly stretched when the carrier is pulled away from the applique. This irreversible stretching can harm the applique's retroreflective performance.\nA second drawback is that when the applique is secured to the substrate with heat, residual carrier web material can remain on the exposed surface of the microspheres, causing a reduction in retroreflective performance. Also, the thermoplastic carrier limits the quantity of heat that can be applied to the applique because, if too much heat is applied to the carrier, the carrier can stick to the microspheres and can become very difficult to remove. Limiting the heat supplied to the applique can cause a poor bond between the applique and the substrate.\nAnother drawback is that known transfer sheets are supplied to the garment assembler with a release liner that, like the carrier web, also becomes discarded as waste. Of course, it is not beneficial from an economic or environmental standpoint to discard the release liner or the carrier web as waste."}
{"text": "A wireless communication device, for example, an Access Point (AP), may implement a Space Division Multiple Access (SDMA) mechanism to transmit a multi-user Multiple Input Multiple Output (MIMO) wireless frame including a plurality of MIMO frames to be simultaneously received by a plurality of wireless communication stations.\nThe multi-user MIMO frame may include a non-beamformed portion and beamformed portion. The non-beamformed portion may include any suitable training symbols and/or signaling fields, which may also be received by other stations. The beamformed portion may be transmitted to the plurality of stations using a respective plurality of beamforming configurations."}
{"text": "Fingerprint image sensors are used in many applications. In one system, a sensor and a host platform together form a fingerprint image processing system. The sensor reads the fingerprint image and converts it to a digital representation. This digital image is transmitted to the host. The host extracts the unique features of the fingerprint, which include minutiae (ridge endings and bifurcations), ridges, valleys, pores, and scars, and converts them into a template, a more compact digital representation of the image. In a first step, the host processes the image and extracts the template, and later compares the template to those stored for authorized users. If a match is found, in the second step the system authenticates the user or performs some other task.\nAs part of sensing a fingerprint image, the sensor continually scans its surface for any image and periodically transmits an interrupt to the host. After receiving the interrupt, the host reads and processes the corresponding digital image data. If the host determines that the image data correspond to a fingerprint, the host can authenticate the user or perform another task. Otherwise, the host ignores the image data and returns to tasks it was performing when it was interrupted. In such systems, smudges or other residual material (films, perspiration, dirt, oils, etc.) are read by the sensor and sent to the host for processing. This requires valuable processing time of the host to determine whether the image is of a usable print or is even a print at all. This use of processing resources is wasteful, unnecessarily generating interrupts to the host processor, diverting the processor from performing other tasks."}
{"text": "Computer graphics systems are widely used in personal computer interfaces, desk top publishing, computer aided design applications, scientific simulations, military and industrial process control and communications networks, and in various education, training and entertainment fields. The breadth of utility of computer graphics systems in so many diverse areas is based upon the singular effectiveness of communicating large volumes of information by means of visual images and the ease, and often the enjoyment, with which humans accept visually presented data. The ability of the computer industry to consistently provide increasingly efficient and economically produced graphics hardware is, however, dependent to large degree upon the costs of electronic components and the development of improved system architectures and methods for graphics image generation.\nMany conventional interactive computer graphics systems employ an applications object model library, an applications program to access, create and process the library of object models, and an image generation system to convert the object models to screen images which may be displayed on a display device such as a computer monitor. To increase usefulness with current display technologies, it is beneficial for an image generation system to be capable of producing a two dimensional frame or screen defined in screen coordinates from a selected group of objects defined in model space. However, such a conversion from an object model to a format suitable for display as a screen is a computationally intense task which can require high levels of a computing systems resources. Thus, the techniques employed to create the contents of a frame (i.e. a complete screen image) and the computer hardware used to generate the data required to describe a frame in many cases determine the effectiveness of a computer graphics system. This relationship is especially true in the area of high performance and intensively user interactive computer graphics systems.\nIn a conventional object model oriented graphics system, the object models, or more simply objects, are digitally formatted geometric descriptions of two or three dimensional objects. An object's description typically consists of a set of geometric primitives. A primitive can be a point, line or polygon defined by one or a set of vertices defined in model space coordinates. The geometric primitives combine to create an object model.\nIn an interactive graphics system, the applications program responds to user input by selecting objects from the applications model library. In fact, certain graphics systems allow the applications program to wholly create new object models. The application program informs the graphics system as to which models are to be displayed and how they are to be displayed, along with the geometric description of each object on a frame by frame or screen by screen basis.\nTo create a screen image the application program defines the description of the image in terms of the object models in the image. Modeling transformation and the viewing operation (i.e. converting from a 3D model to a two dimensional view of that model) are next sequentially performed, followed by rasterization. Modeling transformation and the viewing operation are referred to as front-end or geometry processing. Rasterization, or back-end processing, includes the steps of visible surface determination, scan conversion and shading/illumination.\nEarly raster display architectures provided a system consisting of a system communications bus, a central processing unit(CPU), system memory, a frame buffer, a video controller and a video monitor. The frame buffer could consist of a dedicated memory device or devices, or it may have resided anywhere within the system memory. The CPU, system memory and video controller were all three connected to the system communications bus, through which data, control and status signals were transmitted. In such a system the CPU was required to carry out all of the functions in generating a screen image. This system was not efficient because a general purpose CPU was not well suited to the highly iterative tasks of generating a screen image.\nAn alternate architecture further included a separate display processor and display processor memory through which the video controller was connected to the system communications bus. The display processor performed many of the calculations required to generate frame data, and, therefore, reduced the computational load placed on the system CPU by the applications program. The Texas Instruments TMS34020 peripheral display processor is an example of a device intended for use as a display processor. However, simply segregating the display functions from the CPU in certain instances may not effectively support graphic intense applications such as full motion video and multimedia entertainment applications.\nThe application of multi-processing methodologies, to include pipelining and parallelized computing, has also been attempted to increase the speed with which the large volumes of computation required in graphics image computations could be performed. The inclusion of pipeline and parallel processors are fundamental to the performance of many high performance graphics systems such as the Pixel-Planes and Pixel Flow systems described in U.S. Pat. Nos. 5,388,206 and 5,481,669, to cite examples. These systems include a plurality of logic enhanced memory elements. Each logic enhanced memory element possesses a dedicated arithmetic logic unit (ALU) and a small strip of digital memory. The Pixel-Planes design further includes a linear expression tree, which evaluates linear expressions devised in the form of f(x,y)=Ax+By+C in parallel for every pixel of the screen or region of the screen. Each processing element is dedicated to calculating values for a single pixel. There is only limited capacity for direct interprocessing element communication and the processing elements are, therefore, constrained in the types of functions that they may perform.\nBecause parallel processing has such advantages in image generation, it is desirable to utilize a processing element array such as that used in Pixel-Planes and Pixel-Flow. However, these systems typically have separate geometry processors which handle the conversion of geometric primitives from model coordinates to screen coordinates. Thus, additional hardware or processing capabilities are required or the central processing unit must perform the transformations from model space to screen space.\nDespite the advances in image generation systems, additional work is necessary to reduce the amount of hardware required to provide high speed image generation and to off-load as much processing function from the CPU as is possible."}
{"text": "Thanks to the progress in various technological fields, 3C products (Communication, Consumer electronics, and Computer products) have been constantly improved. For instance, the currently available computer systems now have quicker and quicker computing speed, and the central processing unit (CPU) or the microprocessor of a computer system performing the major computing function thereof has more powerful processing capability. However, the size of the CPU or microprocessor for the computer system does not increase along with the enhanced ability thereof, but is even reduced in some cases. That is, the density of the integrated circuits provided on the CPU or the microprocessor is obviously increased. Thus, heat produced by the CPU or the microprocessor during operation thereof also largely increases to adversely raise the surface temperature of the CPU or the microprocessor even to a level higher than 100° C.\nTo prevent the CPU or the microprocessor from being damaged or causing burnout of nearby elements due to the high temperature thereof, it is necessary to have some mechanism for timely removing the produced heat from the CPU or the microprocessor. Generally, such heat removing mechanism includes a thermal module having a base attached to a surface of a heat-producing element, such as the CPU, the microprocessor, a graphics chip, or a south and a north bridge chip, for absorbing the produced heat; and a heat pipe transferring the heat from the base to a plurality of radiating fins, which is extended through by the heat pipe, so that the produced heat is finally radiated from the radiating fins into ambient air. For the thermal module to stably mount to a mainboard of the computer system, a fastening device is usually needed to connect to the radiating fins extended through by the heat pipe, and a fastening element, such as a screw, is then downward extended through the fastening device to screw into the mainboard.\nFIGS. 1A and 1B as well as FIGS. 1C and 1D showing a first and a second conventional manner, respectively, for fastening a thermal module to a mainboard. The thermal module includes a base 10, a heat pipe 11, a plurality of radiating fins 12, and a fastening device. The heat pipe 11 has an end attached to one side of the base 10, and another side of the base 10 is in contact with a heat-producing element, such as a CPU, a microprocessor, a graphics chip, a south bridge chip or a north bridge chip (not shown). Another end of the heat pipe 11 is extended through the radiating fins 12. The radiating fins 12 has a recess 121 formed on one side facing toward the fastening device 13 for receiving the fastening device 13 therein.\nThe fastening device 13 is made of a metal material, such as iron or copper, and has an end welded to a predetermined position on the radiating fins 12, so that a weld layer 15 is formed between one end of the fastening device 13 and the radiating fins 12 to permanently connect the fastening device 13 to the radiating fins 12. Another opposite end of the fastening device 13 is provided with a fastening hole 131, via which a fastening element 17 is downward extended to screw into a mainboard (not shown), so that the radiating fins 12 are firmly attached to the mainboard via the fastening device 13.\nWhile the fastening device 13 enables fixed connection of the radiating fins 12 to the mainboard, the radiating fins 12 and the fastening device 13 welded thereto could not be separated from one another for reworking unless the weld layer 15 is completely removed from between the radiating fins 12 and the fastening device 13. However, when trying to remove the weld layer 15, it is almost inevitable to cause damages to the fastening device 13 and/or the radiating fins 12.\nFurther, when using the fastening device 13 to fasten the radiating fins 12 to the mainboard, the fastening device 13 is fixedly connected to the radiating fins 12 via the weld layer 15 and could not be freely adjusted to different positions to adapt to different designs. As a result, the conventional fastening device 13 has low applicability and is inconvenient for use and must be produced in different dimensions to match different design requirements, which inevitably increases the overall cost for mounting the thermal module to the mainboard.\nIn brief, the conventional fastening device for thermal module has the following disadvantages: (1) not allowing easy reworking; (2) increasing the cost for mounting the thermal module; and (3) not convenient for use.\nIt is therefore tried by the inventor to develop an improved fastening structure for thermal module, so as to eliminate the above drawbacks."}
{"text": "This invention relates to a low profile catheters such as balloon catheters used to increase the lumen size of a blood vessel, and to guide wires used with such catheters.\nCatheters such as balloon catheters have been used for some time to increase the lumen size of vessels such as coronary blood vessels. Typically, a flexible guide wire is inserted into a selected vessel, and the guide wire then guides the balloon catheter to the treatment site. Packard U.S. Pat. No. 4,646,742 describes one prior art balloon catheter used in an angioplasty procedure. A coil spring guide wire passes through the catheter, beyond the distal end of the catheter and the balloon. Fogarty U.S. Pat. No. 4,606,347 and Samson U.S. Pat. No. 4,641,654 (FIG. 1b) show other balloon catheters which utilize such guide wires. In the catheter shown in the Fogarty patent a seal is provided near the distal end of the catheter between the catheter and the guide wire.\nThis prior art approach requires that the guide wire itself be sufficiently rotationally rigid to transmit rotation from its proximal to distal ends when it is rotated to steer the distal end of the guide wire into a desired vessel. As catheters and guide wires are designed for use with increasingly smaller diameter blood vessels, it may be difficult to ensure that the wire has an adequately small diameter yet is adequately torsionally rigid.\nAnother approach of the prior art is illustrated in the systems described in Mar U.S. Pat. No. 4,771,778, Mar U.S. Pat. No. 4,793,350, Samson U.S. Pat. No. 4,582,181, and Samson U.S. Pat. No. 4,641,654 (FIG. 1a). In these catheters the guide wire is secured to the distal end of the catheter, distally of the balloon, such that a fixed length of the guide wire extends beyond the distal end of the catheter. With this approach, the guide wire is formed as a permanent part of the catheter, and it is therefore not possible to remove the catheter from the blood vessel without simultaneously removing the guide wire. In certain medical procedures it is preferred to leave the guide wire in place at the treatment site while a first catheter is removed from the guide wire. Then a second catheter, possibly different in shape or function from the first, can be guided by the guide wire back to the original treatment site quickly and reliably. Of course, such procedures cannot be performed with a catheter in which the guide wire is permanently affixed to the catheter.\nHorzewski U.S. Pat. No. 4,932,959 discloses a vascular dilation catheter which receives a removable guide wire. The catheter shaft defines an annular portion which may be inflated to releasably secure the catheter to the guide wire to allow the two to be advanced as a single unit. The torque transmitting characteristics of the annular inflatable portion are not discussed. However, it is clear that the inflatable portion must be inflated from the distal end of the catheter to secure the catheter to the guide wire and then deflated from the distal end of the catheter to release the catheter from the guide wire. The time required to inflate and deflate the annular portion may limit the speed with which the guide wire can be shifted between the secured and released modes of use.\nAccordingly, it is an object of this invention to provide an improved low profile catheter for increasing the lumen size of a blood vessel, which utilizes the catheter shaft to transmit torque to the guide wire at a point near the distal end of the guide wire.\nIt is another object of this invention to provide such a catheter in which torque is transmitted from the catheter shaft to the guide wire while permitting relative axial movement therebetween.\nIt is another object of this invention to provide such a catheter in which torque is automatically transmitted from the catheter to the guide wire when the catheter is rotated with respect to the guide wire, and wherein the guide wire is automatically permitted to move longitudinally when the guide wire is advanced with respect to the catheter.\nIt is yet another object of this invention to provide such a catheter in which the guide wire is shaped to allow the catheter to be removed from the guide wire, in order to allow the guide wire to be left in place at the treatment site while the first catheter is removed and replaced with a second."}
{"text": "After “cloud computing” emerges, more and more service providers release cloud computing services, and more and more users use cloud computing services and enjoy core benefits brought by cloud computing: low-cost construction business.\nA cloud storage service is a cloud computing mode that uses storage space as resources and allows a user to use the storage space as required, without worrying the depletion of storage space on an information technology (Information Technology, IT) device. The appearance of a large number of cloud storage services makes it possible to construct a low-cost file storage system.\nConstructing a file storage system directly on the basis of the capability provided by a shared cloud storage service is the simplest scheme, but for an enterprise, the top concern is security and reliability of data assets. However, facts prove that the actual performance of each cloud storage service falls short of users' expectations. Therefore, private cloud storage is currently a mainstream scheme to construct a file storage system by using the cloud storage technology.\nPrivate cloud storage, built upon an enterprise's infrastructure, is open only to internal users of the enterprise. Therefore, the enterprise has the control right of the storage system and can reuse the enterprise's existing storage resources to a certain extent. Private cloud storage is a complex system formed by multiple portions including networks, storage devices, servers, application software, public access interfaces, and client programs. The portions use storage devices as the core and externally provide data storage and a business access service through application software. The construction of such a private cloud storage system requires an enterprise to provide sufficient storage devices that can be virtualized and sufficient sites, to set up a network for connecting the storage devices, and then to deploy, in one by one mode, templates to run at the infrastructure layer, management layer, and application interface layer on the network. Therefore, it can be seen that the construction of a private cloud storage system is complex and implementation is relatively difficult."}
{"text": "This invention relates to a single frequency multi-transmitter telemetry system in which the transmitted data is retained when only one transmitter is transmitting and is discarded when more than one transmitter is transmitting.\nResearch programs in general aviation require the instrumentation of aircraft for test flight purposes. Space, weight, and power consumption are critical factors to be considered in the design of the system and instrumentation of the aircraft. Installation rapidness is also a design consideration in order to make the system cost-effective; primarily because of manpower and aircraft leasing costs.\nIt is the primary object of this invention to provide a telemetry system that will improve the test flight data availability, accelerate the installation process, and reduce data cost.\nAnother object of this invention is to provide a simple single frequency telemetry system that will deliver a substantial amount of data at low cost.\nOther objects and advantages of this invention will be apparent hereinafter in the specification and drawings."}
{"text": "The present invention relates to a drilling apparatus for drilling holes in structures and, in particular, to a drilling apparatus that drills holes in multi-layer structures, and that is adjustable to ensure the drilling tool is maintained in a predefined relationship, such as by being normal, to the area of the structure to be drilled, regardless of the shape of the surface of the structure.\nThe fasteners that hold together multi-layer structures, particularly those structures that are subject to significant dynamic forces and/or pressure over their lifetime, such as aircraft and other vehicle bodies, bridges, buildings, and others, must be properly secured in order to ensure that the structure will perform as intended over its lifetime. In order to properly secure a fastener in a multi-layer structure, the fastener hole must not have any sharp edges, i.e., burrs, the structure must be free of debris between the layers, and any sealant applied between the layers in order to make the structure air and/or water tight must be sufficiently squeezed out. When excess sealant is present between the layers, the distance between the layers is increased and/or uneven, which may be referred to as a “gasket” condition. Thus, if burrs, debris and/or excessive sealant are present, then the layers cannot be properly fastened and the layers may suffer corrosion, cracking and/or premature fatigue failure, which generally renders the structure ineffective for its intended purpose and, therefore, subject to the expense of repair or replacement.\nThus, ensuring that a burr-less hole is drilled, that there is no debris between the layers, and that sealant is properly applied between the layers is an integral part of clamping multi-layer structures together. In the aerospace industry, for example, a significant amount of time and labor is expended ensuring that the holes through the various layers of the aircraft structure are appropriately drilled, cleaned, sealed and fastened. Initially, the layers of materials that form the structure are loosely assembled without sealant, and drill templates are aligned and attached to the structure in the areas to be drilled. A drill operator, guided by the drill template, then drills holes through the layers of materials typically using a manual drill motor. As the hole is being drilled through the layers, the drill bit tip pushes with the full feed force applied to the drill motor. This can cause a gap to develop between a layer currently being drilled and the underlying layer, particularly when the layers are a stack-up of thin material. The gap between the layers causes burrs about the hole and debris is likely to gather between the layers. Thus, once the holes are drilled, the layers must be disassembled, the burrs must be removed from about the holes, and the debris must be cleaned from the surfaces of the layers, all of which is a time-consuming and labor intensive process.\nSealant is then applied to the layers prior to re-assembling the layers. In order to ensure the layers are properly sealed to provide an air and water-tight seal, a generous amount of sealant is applied to the layers. Clamps that extend through the holes, such as KWIK-LOK™ clamps commercially available from Zephyr Manufacturing Company, Inglewood, Calif., must be placed through each hole of the reassembled layers in order to squeeze out the sealant to prevent excessive “gasket” between the layers before the sealant dries. The extra sealant squeezes out around the clamps and must be cleaned from the structure and the clamps during clamping and/or after the clamps are removed.\nOnce the sealant is cured, the clamps may be removed and the holes may be countersunk. In order to countersink a hole, a countersink drill bit and microstop countersink cage are attached to the drill motor and the operator revisits each hole to drill the countersink. The holes are then inspected to ensure they were properly countersunk. The holes may be inspected by checking the countersink of each hole with a measuring tool, or by installing the fastener to check if it fits properly within the hole and countersink. If the holes are satisfactory, then fasteners may be installed and fastened with nuts or swage lock collars. Overall, this process is expensive, laborious, and time-consuming. In addition, the integrity of the resulting holes depends upon the completion of many manual processes, which creates a risk that certain steps may be performed inadequately or completely overlooked.\nFor structures that have a flat outer surface, a positive, power feed drill motor with a collet mandrel is capable of drilling, reaming and countersinking holes in multi-layer structures in one step. Examples of such positive power feed drill motors with a collet mandrel are the Q-matic/Space-matic drill motors, and, in particular the Q-matic 15SC Self-Colleting Drill Motor, produced under the name Quackenbush, and commercially available from Cooper Tools, Inc. To utilize these drill motors, a hole is initially drilled through the layers of material by conventional means, such as by using a conventional drill motor to drill a hole through the layers, then de-burring the hole and cleaning the debris from between the layers. Once the structure is reassembled, the collet is extended through the hole and secured with the mandrel on the side of the structure that is opposite the drill motor. The collet mandrel is typically attached to a platform of the drill motor such that the platform sets on the surface of the structure. The platform also generally has an opening that is aligned with the structure to be drilled, and the drilling tool of the drill motor may extend through the platform opening to drill through the structure. Thus, the collet mandrel securely holds the layers of the structure together such that holes may be drilled and countersunk within a predetermined distance from the collet mandrel in one step without the risk of debris accumulating between the layers. Once all of the desired holes have been drilled within the predetermined distance from the collet mandrel, the collet mandrel is removed from the hole and moved to one of the holes recently drilled or a new hole must be drilled for the collet mandrel, prior to repeating the process described above to form additional holes.\nHoles are preferably drilled at a predefined orientation, such as a normal orientation, relative to the surface of the structure. As long as the surface of the structure is flat, the drilling tool of the drill motor will typically be normal to the structure to be drilled because the platform lies flat on the surface. Thus, so long as the platform is normal to the drilling tool, then the surface of the structure will also be normal to the drilling tool. If, however, the surface is contoured or three-dimensional, the drilling tool may not be normal to the area of the structure to be drilled because the platform cannot lie flat on the surface. Thus, the drill motor with the collet mandrel described above may include an adjustable foot pad to adjust the position of the platform, and therefore the drilling tool until the drilling tool is normal to the area of the structure to be drilled. The adjustable foot pad is typically attached to the drill motor platform opposite the side of the platform with the opening, and is adjustable to various positions based upon the radius of curvature of the surface and the distance from the foot pad to the drilling tool. This design is suitable for simple contours with only one radius of curvature between the foot pad and the drilling tool, i.e., for surfaces that curve in only one dimension, because the adjustable foot pad maintains contact with the contoured surface across the full width of the pad. If, however, more than one radius of curvature is present between the foot pad and the drilling tool, such as in complexly curved surfaces exhibiting curvature in two or more dimensions, the pad cannot maintain contact with the surface across the width of the pad because the curvature of the surface may be different at one side of the foot pad than it is at the other side of the foot pad. Thus, the conventional adjustable foot pad is not capable of ensuring that the drilling tool is normal or in any other predefined relationship to the surface of the structure to be drilled when the surface of the structure is complexly curved. As such, to drill a hole in a multi-layer structure having a complex curvature, the holes must be drilled in the conventional manner of drilling, de-burring, cleaning the layers, then reassembling and sealing the layers before countersinking the holes, which is a time-consuming and labor-intensive process, as discussed above.\nTherefore, a need exists to drill holes in multi-layer structures having surfaces that are at least partially complexly curved in a more efficient manner, such as by reducing or eliminating the time and expense involved in disassembling the drilled layers, de-burring the holes, cleaning the debris from the layers, reassembling the layers, clamping the holes, waiting for the sealant to dry, and revisiting the holes to drill the countersink. In particular, there is a need to efficiently drill holes in multilayer structures having surfaces that are at least partially complexly curved in one operation, such that the time and money involved in the manufacture of such multilayer structures is significantly reduced."}
{"text": "There are prior art speed variators which have a driving wheel that engages, by rolling, on the outer surface of a conical bell-shaped body (driven unit) in such a way that it can slide along the axis of the bell-shaped body to vary the radius of the circumference traveled by the wheel on the bell-shaped body, thus varying the gear ratio. In these variators, the axes of rotation of the bell-shaped body and the wheel are parallel.\nDisadvantageously, in conventional variators of the type just described, in the contact area between the wheel and the bell-shaped body, micro-slippages are created which result in a dissipation of power, penalising the overall performance of the reduction stage. Such slippages are mostly caused by a kinematic coupling which is not optimum between the wheel and the bell-shaped body.\nMoreover, the local deformability of the wheel and/or the bell-shaped body generates hysteresis cycles which have the effect of dissipating a further fraction of the power transmittable."}
{"text": "The present invention relates generally to recognition of user inputs to a computer system such as a pen-based computer system and, more particularly, to methods for analyzing the user inputs and providing interpretations of those user inputs.\nA pen-based computer system is a small, often hand-held, computer system where the primary method for inputting data includes a \"pen\" or stylus. A typical pen-based computer system is housed in a small, rectangular enclosure, and has a dual-function display assembly having a viewing screen exposed along one of the planar sides of the enclosure. The dual-function display serves as both an input device and an output device. When operating as an input device, the display assembly senses the position of the tip of the stylus on the viewing screen and provides this information to the computer's central processing unit (CPU). Some display assemblies can also sense the pressure of the stylus on the screen to provide further information to the CPU. When operating as an output device, the screen displays computer-generated images developed by the CPU.\nThe dual-function display assemblies of pen-based computer systems permit users to operate the computer as a computerized notepad. For example, graphical images can be input into the pen-based computer by merely moving the stylus over the surface of the screen. As the CPU senses the position and movement of the stylus, it generates a corresponding image on the screen to create the illusion that the stylus is drawing the image directly upon the screen. With suitable recognition software, text and numeric information can also be entered into the pen-based computer system in a similar fashion.\nWhat is needed is a system for recognition of user inputs to a computer system such as a pen-based system and for analysing the user inputs to determine possible meanings."}
{"text": "In a wireless communication networks, there are communication nodes, for example base stations, where base station products are provided by different suppliers. Today there is typically one base station product per frequency band. In the future the number of frequency bands is expected to grow rapidly, which would require a lot of different base station products supported by each supplier. Each base station product is expensive to develop and to service, which will make it very costly to have a plurality of different base station products.\nOne way for reducing the number of base station products is to design multiband base stations, which means that one base station product could be used for many different frequency bands. One essential part of a multiband base station is multiband filters, which is used to filter the signals before they are transmitted/received. However, multiband filters can only be designed for relatively low signal power levels, which makes it difficult to use them for transmission in macro and micro base stations.\nIn base stations, there is usually one filter per transmitter radio arrangement, where a transmitter radio arrangement comprises one or more power amplifiers. By increasing the number of transmitter radios arrangements, it is possible to reduce the output power per radio and hence the power going through the corresponding filters. With large enough number of transmitter radio arrangements, it is then possible to reduce the power per filter such that multiband filters are enabled. However, such an arrangement would be rather costly due to the larger number of transmitter radio arrangements.\nThere is thus a need for a multi-band communication node where the number of transmitter radio arrangements is not increased compared to a single-band communication node."}
{"text": "Field of the Invention\nThe present invention relates to a burner configuration for a firing installation, in particular a gas turbine combustion chamber, having a main burner and, disposed centrally within it, a primary pilot burner which is used for igniting and/or for stabilizing the combustion of the main burner. In addition, the present invention relates to a method of operating such a burner configuration.\nBurner configurations of the generic type are often configured for different fuels and different modes of operation and, in view of the continually more strict regulations worldwide on the emission of pollutants, great efforts are being made to improve the configurations and modes of operation in such a way that the emission of pollutants is reduced.\nParticularly environmentally friendly burners with a low emission of NOx are known from European Patent EP 0 108 361 B1 and European Patent EP 0 193 838 B1. Because the present invention deals precisely with a development of such burners, reference is made to the complete content of these two publications.\nPublished, Non-Prosecuted German Patent Application DE 196 10 930 A1 discloses a burner for a heat generator. In the burner, which consists essentially of a swirl generator for a combustion airflow and of a device for injecting a fuel into the combustion airflow, a mixing section is disposed downstream of the swirl generator mentioned. Within a first partial section, the mixing section has a number of transition ducts extending in the flow direction, which transition ducts ensure the continuous transfer of the flow formed in the swirl generator into a downstream mixing tube. The outlet plane of the mixing tube relative to the combustion chamber is configured with a final edge which serves to stabilize and increase a reverse-flow zone forming in the flow. A number of mixing elements, which serve to form a mixture of combustion air and a fuel, are provided concentrically with the mixing section. This mixture of the respective mixing elements then forms one pilot stage of the combustion chamber.\nAlthough, in the known burner configurations, critical conditions of the burner can be substantially avoided by design measures and the mode of operation, there is nevertheless the possibilityxe2x80x94particularly when the burner configuration is operated in premixed operation, i.e. when fuel and air are already intensively mixed with one another before the combustion zonexe2x80x94that the combustion in a firing installation exhibits fluctuations and sometimes becomes unstable, which can, for example be expressed by a firing installation rumble. The flame can then also, on occasion, flash back into a partial region of one of the burners, which is undesirable because of the increases in temperature of burner components caused by it and because of an increased emission of pollutants which then occurs.\nIt is accordingly an object of the invention to provide a burner configuration with primary and secondary pilot burners which overcomes the above-mentioned disadvantages of the prior art devices of this general type, which creates increased stabilization of the combustion and, helps avoid firing installation rumble due to unstable combustion and flashback of the flame into partial regions of the burner.\nWith the foregoing and other objects in view there is provided, in accordance with the invention, a burner configuration for a firing installation such as a gas turbine combustion chamber. The burner configuration contains a main burner having an outlet region and a primary pilot burner disposed centrally within the main burner. The primary pilot burner ignites and stabilizes a combustion of the main burner. A fuel supply line and a secondary pilot burner having a peripheral annular duct with a plurality of outlet openings formed therein and disposed in the outlet region of the main burner, are provided. The peripheral annular duct is connected to the fuel supply line and is in connection with the outlet openings, the outlet openings being directed in a direction towards a center of the outlet region.\nSuch a burner configuration for a firing installation, in particular a gas turbine combustion chamber, has the main burner and, disposed centrally within it, a primary pilot burner which is used for igniting and/or for stabilizing the combustion of the main burner. Details of such burners, which are also called hybrid burners because of the different possibilities with respect to mode of operation and fuel supply, are described, for example, in European Patent EP 0 108 361 B1 and European Patent EP 0 193 838 B1. In the invention of the instant application, such a burner configuration is now additionally equipped with a secondary pilot burner, the latter exhibiting a plurality of outlet openings in the outlet region of the main burner. Because fuel, or a mixture of air and fuel, emerges from these outlet openings, a ring of flame is formed which surrounds the flame emerging from the main burner approximately as an annulus and, by this, additionally stabilizes the flame. Whereas the primary pilot burner supports the combustion centrally, more or less from the inside out, the secondary pilot burner can contribute to the stabilization from the outside without negatively affecting the principle of combustion with recirculation, which is possessed by the burners of the prior art. The outer recirculation region, which is important for stable combustion in the burner flame, is in fact strengthened by the secondary pilot burner.\nA particularly preferred configuration is for the secondary pilot burner to be formed by a peripheral annular duct in the outlet region of the main burner, which annular duct is connected to a fuel supply line and is in connection with the outlet openings. The peripheral annular duct surrounds the outlet region of the main burner more or less like a collar in which the outlet openings are disposed, preferably directed toward the inside in the direction of the center of the outlet region.\nTo avoid any limitation to the number or the cross-sectional area of the outlet openings, the peripheral annular duct is subdivided, in a particularly preferred embodiment example, into a lower annular duct and an upper annular duct, which are connected together by a throttle location, the upper annular duct being connected to the fuel supply line and the lower annular duct being in connection with the outlet openings. In this way, the fuel quantity for the secondary pilot burner is determined by the throttle location and not by the cross section of the outlet openings and their number. In this way, the lower annular duct acts as a resonance chamber so that the fuel quantity emerging at the outlet openings is influenced by the external pressure which, given appropriate design and dimensioning of this lower annular duct, acts in a strongly damping manner in the case of pressure vibrations, i.e. during burner configuration rumble for example.\nAccording to the invention, the outlet openings can preferably be disposed in a ring, at approximately equal distances apart, in the outlet region of the main burner. This configuration is symmetrical to all sides and leads to uniform stabilization on all sides.\nAs an alternative, it is also possible to dispose the outlet openings in a ring, at different distances apart, in the outlet region of the main burner, this being useful, particularly in the case of combustion chamber configurations with a plurality of burners, in order to avoid accumulations of outlet openings in regions in which two burners come very close together. In addition, a non-uniform distribution of the outlet openings has a smaller tendency to vibration, and therefore to instabilities, than a uniform distribution because, in the case of pressure shocks, different pressure wave traverse times to or from the individual outlet nozzles occur due to different distances and this leads, at least partially, to mutual extinguishing of such pressure waves and therefore damps instabilities.\nThere is also a large degree of freedom with respect to the shape of the outlet openings. Particularly preferred, according to the invention, are round or oval outlet openings, which are particularly easy to manufacture.\nDepending on the burner geometry, however, the outlet openings can have a polygonal edge, in particular in the form of a triangle with an apex pointing in the gas flow direction. The burner configuration with the secondary pilot burner is configured, according to the invention, in such a way that between 0 and 10% of the fuel quantity required at base load of the burner configuration can be supplied to the outlet openings, in particular in the form of fuel gas. Depending on the mode of operation of the burner configurations, it can be useful to appropriately throttle the fuel supply to the central primary pilot burner, which is now supported, according to the invention, by the secondary pilot burner, in order not to unduly increase the proportion of fuel which is supplied by the pilot burners overall.\nIn the burner configurations of the prior art, it is known art for the main burner and the primary pilot burner to be operated in premixed operation, in diffusion operation or in a mixed operation of both types of operation, it being also possible to configure the main burner and the primary pilot burner for optional operation with liquid and gaseous fuels. All these modes of operation are also possible in the case of a burner configuration according to the present invention. In addition, it is also possible to operate the secondary pilot burner as a premixing burner with a premixed mixture of fuel and air, which particularly contributes to reducing the generation of oxides of nitrogen.\nThe invention also relates to a method of operating the burner configuration, in particular to avoid instabilities in the combustion chamber, the associated rumble and the possible flashback of the flame in partial regions of the burner, the burner configuration having the main burner stabilized by the primary pilot burner and, according to the invention, the flame of the main burner being additionally stabilized by the secondary pilot burner, the secondary pilot burner having a plurality of outlet openings distributed around the outlet region of the main burner, which outlet openings are supplied separately, using open-chain control or closed-loop control, with fuel or fuel mixture. This method according to the invention provides a further parameter for the particularly reliable and environmentally friendly control of a burner configuration. The stabilization of the flame of the main burner by the outer secondary pilot burner inhibits unstable operating conditions such as rumble in the firing installation and, therefore, also inhibits flashback of the flame in partial regions of a burner.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in a burner configuration with primary and secondary pilot burners, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawing."}
{"text": "1. Field of the Invention\nThe present invention relates to a coating device and a coating method. More specifically, it relates to a coating device and a coating method capable of performing a stable and even coating operation, even in the case of coating a coating liquid agent on the surface of band-shaped substrate, with the substrate running at a high speed.\nFurthermore, the invention relates to a coating device and a coating method using a bar coater. More specifically, it relates to a bar coating device and a bar coating method capable of coating a coating liquid agent on an object to be coated being conveyed in a constant direction at a high speed, preferably adopted for cases where the object to be coated has a rough surface.\nFurthermore, the invention relates to a bar coating device comprising a bar and a primary side weir member disposed upstream of the bar, capable of effectively preventing drying or adhesion of a coating liquid on the primary side weir member, and a coating method using the coating device.\n2. Description of the Related Art\nPhotographic sensitive materials, printing papers, magnetic recording materials, coated metal plates, planographic printing plate precursors, or the like have been produced by coating a coating liquid such as a plate-making layer forming liquid on a substrate such as a support web.\nA bar coater, a slide bead coater, an extrusion coater, of the like have been used for the application of the coating liquid. In particular, the bar coater has been used widely because it is easy to handle.\nThe bar coater comprising a bar that is rotatable in the same or opposite direction relative to the conveyance direction of the substrate while contacting the lower surface of the continuously running substrate has been conventionally used. In that bar coater, by ejecting the coating liquid upstream of the bar in the substrate conveyance direction (hereinafter referred to simply as “upstream”) while the substrate is running so as to form a bead and taking up the coating liquid by the rotation of the bar, the coating liquid is adhered onto the lower surface of the substrate. The bar may be slaved to the running of the substrate and rotated thereby.\nMoreover, in addition thereto, a bar coater disclosed in the specification of the Japanese Utility Model Application Laid-Open No. 63-126213, comprising a first weir plate disposed adjacent to a bar and upstream of the bar, with the thickness of the upper end portion thereof tapered toward the downstream conveyance direction of the substrate hereinafter simply referred to as “downstream”) and the upper end portion bent toward the bar, having a flat surface of the length 0.1 to 1 at the top portion; and a bar coater disclosed in the official gazette of the Japanese Patent Application Publication (JP-B) No, 58-004589 comprising a bar, a first weir plate disposed on the upstream side thereof, with the upper end portion thickness made thinner toward the downstream side, and a second weir plate disposed downstream of the bar, have been commonly used.\nFor the above bar coaters, a coating liquid is generally discharged between the first weir plate and the bar so as to be coated on the substrate.\nMoreover, in the case where the weir plate is provided upstream of the bar, those methods of coating the coating liquid on the substrate by ejecting the coating liquid between the bar and the weir member while forcibly rotating the bar at a circumferential speed different from the conveyance speed of the substrate, have mainly been used.\nHowever, for any of the above bar coaters, when the substrate is running at a high speed, stripe-like defects of an equal pitch can easily be generated. Moreover, when the support web is run at a high speed, an accompanying film of air following the support web, can form on the surface of the support web.\nWhen an entrained air film is formed on the surface of the support web, the entrained air film is brought by the support web to the bead as the coating part, and this generates such defects as coating film cuts and coating irregularities in the coating film of the coating liquid formed on the surface of the support web. Thus there is a problem involved in that the coating liquid cannot be coated on in a stable manner.\nMoreover, for the bar coater provided with the weir plate, the coating liquid is intentionally discharged between the first weir plate and the bar, and it cannot always be provided evenly along the substrate width direction.\nTherefore, the coating liquid discharged intermittently dries and adheres on the top portion of the first weir plate so that due to the adhered coating liquid solid component, coated surface problems such as coating streaks and solid component adhesion can be generated on the coated surface of the substrate.\nFurthermore, in order to coat a coating liquid on an object to be coated such as a metal plate, or eliminate excessive coating liquid from the object to be coated (that is called, “measurement”), a bar coating device 502 as shown in FIG. 20 can be used.\nFor the bar coating device 502, a columnar (cylindrical) coating bar 506 is provided for a metal plate 504. The columnar coating bar 506 is conveyed at a constant speed in a direction orthogonal to the direction in which the metal plate 504 is conveyed (arrow F1 direction) so as to be contacted with the coated surface (lower surface) of the metal plate 504. The coating bar 506 is rotated at a circumferential speed equal to the conveyance speed of the metal plate 504 by the friction between the metal plate 504 and the coating bar 506. The rotation brings up a coating liquid 508 to provide a bead 510 between a weir member 512 and the metal plate 504. That is, the coating liquid in the bead 510 is coated on the metal plate 504 and the excessive coating liquid is eliminated from the metal plate 504, i.e., measured.\nHowever, for the bar coating device 502, the so-called entrained air enters in at the time of coating so that the bead 510 is not stable. In particular, the conveyance speed of the metal plate 504 has recently been set at high speeds, so that due to the high speeds, the risk of the inability to maintain stability in the bead 510 is increased. If the bead 510 is instable, for example, the coating liquid 508 cannot be coated evenly in the entire width direction of the metal plate 504, so that it is difficult to obtain an evenly coated surface. This is especially true in the case where the surface of the metal plate 504 is coarse, since entrained air can easily be generated, it is further difficult to obtain an evenly coated surface."}
{"text": "The invention relates to vascular repair devices, and in particular intravascular stents, which are adapted to be implanted into a patient's body lumen, such as a blood vessel or coronary artery, to maintain the patency thereof. Stents are particularly useful in the treatment of atherosclerotic stenosis in arteries and blood vessels.\nStents are generally tubular-shaped devices which function to hold open a segment of a blood vessel or other body lumen such as a coronary artery. They also are suitable for use to support and hold back a dissected arterial lining that can occlude the fluid passageway. At present, there are numerous commercial stents being marketed throughout the world. For example, the prior art stents depicted in FIGS. 1-3 have multiplex cylindrical rings connected by one or more undulating links. While some of these stents are flexible and have the appropriate radial rigidity needed to hold open a vessel or artery, there typically is a tradeoff between flexibility and radial strength and the ability to tightly compress or crimp the stent onto a catheter so that it does not move relative to the catheter or dislodge prematurely prior to controlled implantation in a vessel.\nWhat has been needed and heretofore unavailable is a stent which has a high degree of flexibility so that it can be advanced through tortuous passageways and can be readily expanded, and yet have the mechanical strength to hold open the body lumen or artery into which it is implanted and provide adequate vessel wall coverage. In particular, it would be desirable to have a stent with an end ring configuration that aims at reducing the stent-to-shoulder distance as well as delivering therapeutic drug to the peri-stent area while maintaining a pristine stent deployment. The present invention satisfies these and other needs. That is, the stent of the present invention has a high degree of compressibility to secure it on the catheter and provide a low profile and a high degree of flexibility making it possible to advance the stent easily through tortuous arteries, yet the stent has sufficient radial rigidity so that it can hold open an artery or other blood vessel, or tack up a dissected lining and provide adequate vessel wall coverage."}
{"text": "The subject invention provides for pitch, yaw and clockwise and counter clockwise roll control of a missile and more particularly but not by way of limitation to the use of a plurality of control panels which are slanted for controlling clockwise and counter clockwise roll of a missile.\nHeretofore, there have been various types of missile control systems and arrangements such as drag controls, brake flaps and fins for controlling the operation of the missile. These devices and systems are disclosed in the following United States Patents: U.S. Pat. No. 2,793,591 to Jasse, U.S. Pat. No. 2,941,764 to Lee, Jr. et al, U.S. Pat. No. 2,942,545 to Fogal et al, U.S. Pat. No. 3,004,489 to Griffith et al, U.S. Pat. No. 3,114,315 to Trump, U.S. Pat. No. 3,174,430 to Apotheloz, U.S. Pat. No. 3,188,958 to Burke et al, U.S. Pat. No. 3,343,767 to Cafissi, U.S. Pat. No. 3,588,004 to Suter, U.S. Pat. No. 3,622,103 to Meier.\nNone of the above mentioned patents provide the unique features and advantages of the subject invention."}
{"text": "The present invention relates to cotton plants that produce cotton fibers with improved fiber characteristics such as increased fiber length, increased fiber fineness and higher fiber strength, and it also relates to a method for producing cotton fibers from these cotton plants.\nCotton fibers are usually produced by cultivating cotton plants of the genus Gossypium and collecting the cotton fibers from the capsules (cotton bolls) formed on the cotton plants. There are many varieties of cotton plant, from which cotton fibers with different fiber characteristics can be obtained and used for various applications depending on their fiber characteristics. Cotton fibers are characterized by various properties among which fiber length, fiber fineness and fiber strength are particularly important. Many efforts have been made so far to improve the characteristics of cotton fibers. Attempted improvements have been mainly focused on fiber length and fiber fineness. In particular, there has been a great demand for longer and finer cotton fibers. The variety of cotton plant known as Sea Island is famous for desired fiber characteristics; however, this variety of cotton plant exhibits a poor yield of cotton fibers and therefore the price of Sea Island cotton fibers is very high. If highly yielding cotton plants with cotton fiber characteristics equal to or better than those of Sea Island cotton can be produced, it will be a great contribution to industry.\nThe methods for improving the characteristics or yield of cotton fibers can be roughly classified into the following three categories:\n1. Variety Improvement by Cross Breeding\nThis method has been utilized most widely so far. At present, almost all the cultivated varieties of cotton plant are bred by this method. However, much time is needed for this method and the degree of variability is limited, so that one cannot expect remarkable improvements in the cotton fiber characteristics or in the yield of cotton fibers.\n2. Treatment with Plant Hormone\nPlant hormones such as auxin, gibberellin, cytokinin and ethylene have been put to practical use widely for field crops or horticultural products. The influence of plant hormones on the fiber characteristics of cotton plants, particularly on the fiber elongation mechanism, is known by gibberellin, auxin or brassinolide; however, no effect has been fully confirmed yet, and it cannot be said that these plant hormones are effective for practical use.\n3. Variety Improvement by Gene Recombination Technique\nIn recent years, gene recombinant technique has made remarkable progress, and several reports have been made on the successful variety improvement in certain kinds of plants (e.g., cotton, soybean, corn, tomato) by introduction and expression of a particular gene in these plants to confer a desired genetic trait thereon. There have been developed and put to practical use, for example, cotton plants with improved insect resistance by introduction of a gene coding for BT toxin (i.e., insecticidal protein toxin produced by Bacillus thuringiensis) or cotton plants with improved herbicide (Glyphosate) resistance by introduction of a gene coding for 5-enol-pyruvil-shikimic acid 3-phosphate synthetase.\nIf a gene associated with cotton fiber formation and elongation can be introduced into cotton plants and expressed in sufficient quantities, it would become possible to make a remarkable improvement in the characteristics or yield of cotton fibers. Further, the introduction of such a gene in anti-sense for makes it possible to suppress the action of this gene. In other words, it is believed that the characteristics and yield of cotton fibers can be controlled by introducing a gene associated with fiber formation and elongation into cotton plants, followed by large-scale expression or suppression of the gene. The method using such a genetic engineering technique can be expected to find wide applications, for example, more reliable control of fiber elongation and formation as compared with the conventional plant breeding by cross-breeding and screening. For this purpose, it is required that the mechanisms of fiber elongation and formation are elucidated on the genetic level and some genes closely associated with these mechanisms are discovered and then actually expressed and regulated in the cotton fiber tissues to control the mechanisms of fiber elongation and formation.\nAt present, however, the knowledge of fiber elongation and formation in plants from the viewpoint of molecular biology is very limited. Although many studies have been made on the elongation of plant cells, most of the control factors remain unknown and the control mechanisms have not yet been elucidated.\nA cotton fiber is composed of a single cell that has been differentiated from an epidermal cell of the seed coat, and it develops through four stages, i.e., initiation, elongation, secondary cell wall thickening and maturation stages. More specifically, the elongation of a cotton fiber begins with that of an epidermal cell of the ovule just after flowering and the cotton fiber rapidly elongates and then completes elongation in about 25 days after the flowering. After that, the fiber elongation is stopped, and a secondary cell wall is formed and grown through maturation to become a mature cotton fiber.\nSome reports have been made on the isolation of such genes associated with the elongation and formation of cotton fibers. John et al. describe the isolation of E6 gene that is expressed preferentially in the cotton fiber tissues on the 15th and 24th days after flowering (see John, M. E. and Crow, L. J., Proc. Natl. Acad. Sci. USA, 89, 5769-5773 (1992)) or H6 gene coding for a proline-rich protein that actively functions in the formation of secondary cell walls (see John, M. E. and Keller, G., Plant Physiol., 108, 669-676 (1995)). John further examined the effects of E6 gene on the cotton fiber characteristics by introduction of anti-sense E6 gene into cotton plants to reduce the expression level of endogenous E6 RNA (John, M. E., Plant Molecular Biology, 30, 297-306 (1996)). John, however, reported that although the expression level of E6 gene in fiber tissues is reduced, fiber length, fiber strength and fiber fineness are not significantly changed, and he concluded that E6 gene is not critical to the normal development of cotton fibers. Song et al. identified acyl carrier protein (ACP) cDNA from cotton plants by differential display method, which protein is specifically expressed in cotton fiber tissues (Song, P. and Allen, R. D., Biochimika et Biophysica Acta, 1351, 305-312 (1997)). As the gene associated with the cellulose synthesis in cotton fibers, cDNA coding for a catalytic subunit of cellulose synthase was isolated (Pear, J. R., Kawagoe, Y., et al, Proc. Natl. Acad. Sci. USA, 93, 12637-12642 (1996)). Kasukabe, one of the present inventors, and his coworkers have isolated and identified five genes from cotton plants by differential screening method and differential display method, which genes are specifically expressed at the cotton fiber elongation stage (assignees\"\" own U.S. Pat. No. 5,880,100 and U.S. patent applications Ser. Nos. 08/580,545, 08/867,484 and 09/262,653). Some genes associated with the elongation and formation of cotton fibers have already been isolated from cotton plants; however, none have succeeded in modifying the fiber characteristics of cotton plants in practice.\nThe analysis of molecular mechanisms of plant cell wall construction led to the isolation of endoxyloglucan transferase (EXGT) as an enzyme that catalyzes molecular grafting between polysaccharide cross-links in the plant cell wall matrix (Nishitani, K. and Tominaga, R., J. Biol. Chem., 267, 21058-21064 (1992)). Xyloglucans are polysaccharides that cross-link individual cellulose microfibrils and play the main role in the net work structure of the cell wall involved in the cell elongation. For this reason, the transfer of xyloglucan cross-links by endoxyloglucan transferase is considered one of the important processes in the cell elongation. Some genes coding for endoxyloglucan transferase have been isolated from various plants including tomato (Lycopersicon esculentum) and mouse-ear cress (Arabidopsis thaliana) (Arrowsmith, D. A. and de Silva, J., Plant Mol. Biol., 28, 391-403 (1995); Okazawa, K., Sato, Y., et al, J. Biol. Chem., 268, 25364-25368 (1993)), and further from cotton plants (Gossypium spp.) (Shimizu, Y., Aotsuka, S., et al., Plant Cell Physiol., 38(3), 375-378 (1997)). There is, however, no report that endoxyloglucan transferase genes isolated from various plants or from cotton plants are actually introduced into cotton plants and examined for the effects on the cotton fiber characteristics.\nMost of the plant catalase species are localized on the microbodies and they are detoxification enzymes that decompose metabolically produced toxic H2O2 into water and oxygen. Their functions in plants have not yet been fully elucidated, although some reports have been made that catalase is associated with the low temperature tolerance and pathogen resistance of plants (Sanchez-Casas, P. and Klessing, D. F., Plant Physiol., 106, 1675-1679 (1994); Prasad, T. K., Anderson, M. D., et al., Plant Cell, 6, 65-74 (1994)). There is, however, no report on the relationship between catalase and fiber characteristics of cotton plants.\nPeroxidase is an enzyme that catalyzes the following oxidative reaction by hydrogen peroxide: H2O2+AH2xe2x86x922H2O+A\nPeroxidase genes have been isolated from various plants including tomato (Lycopersicon esculentum) and horseradish (Armoracia rusticana) (Roberts, E., Kolattukudy, P. E., Mol. Gen. Genet., 217, 223-232 (1987); Fujiyama, K., Tekemura, H., et al., Eur. J. Biochem., 173, 681-687 (1988)). Peroxidase is one of the cell wall enzymes and it exists through ionic bonding or covalent bonding in cell walls. One of the functions of peroxidase is the formation of lignin in the secondary components of cell walls. Peroxidase is considered to catalyze, in the presence of hydrogen peroxide, the reaction of forming dehydrogenation polymeric products from the lignin components (e.g., ferulic acid, p-coumaric acid) produced through shikimic acid pathway or cinnamic acid pathway (Gross, G. G. et al., Planta, 136, 271 (1977)). Another report describes that peroxidase is associated with the formation of intermolecular cross-links by the reaction of the respective tyrosine residues of two extensin molecules as the structural proteins of cell walls (Lamport, D. T. A., MSU-DOE Plant Research Laboratory 7th Annual Report, p. 65 (1982)). With respect to the relationship between peroxidase and cotton fibers, Rao et al. examined the peroxidase activity during the development of cotton fibers and showed that the peroxidase activity is lower in the cotton fiber elongation but significantly higher in the secondary wall thickening (Rao, N. R., Naithani, S. C., et al., Z. Pflanzenphysiol. Bd., 106, 157-165 (1982)). These results suggest a possibility that peroxidase is involved in the rigidifying of cotton fiber cell walls. John et al. confirmed that fiber strength can be increased by introduction of a peroxidase gene into cotton plants and over-expression of the gene in the cotton fiber tissues (WO95/08914); however, they failed to obtain significant results on the fiber length or micronaire (i.e., fiber fineness).\nAs described above, some genes associated with the elongation and formation of cotton fibers have already been isolated; however, it seems to be the most important to actually introduce these genes into cotton plants to make sure of their effects on the level of practical use.\nUnder these circumstances, the present inventors have extensively studied the mechanisms of fiber elongation and formation in cotton plants from the viewpoints of molecular biology to improve the characteristics of cotton fibers. As a result, they have found that this purpose can be attained by introducing a gene coding for endoxyloglucan transferase, which is deeply associated with the cell elongation and greatly expressed in the cotton fibers and ovule tissues at the cotton fiber elongation stage, or a gene coding for catalase or peroxidase, which is a hydrogen peroxide eliminating enzyme, into cotton plants and over-expressing these genes in the cotton fiber cells, thereby completing the present invention.\nThus, the present invention provides a cotton plant of the genus Gossypium with improved cotton fiber characteristics, comprising an expression cassette containing a gene coding for an enzyme selected from the group consisting of endoxyloglucan transferase, catalase and peroxidase so that the gene is expressed in cotton fiber cells to improve the cotton fiber characteristics; and cotton fibers and cotton seeds obtained from the cotton plant.\nThe present invention further provides a method for producing cotton fibers with improved fiber characteristics, and cotton fibers obtained by this method. The method includes the steps of: transforming into a cotton plant of the genus Gossypium, an expression cassette containing a gene coding for an enzyme selected from the group consisting of endoxyloglucan transferase, catalase and peroxidase so that the gene is expressed in cotton fiber cells to improve the cotton fiber characteristics; growing the cotton plant transformant to form cotton bolls; and collecting cotton fibers from the cotton bolls of the plant."}
{"text": "1. Field of the Invention\nThe present invention relates to combustion control processes for reducing harmful emissions produced in conventional diesel internal combustion engines.\n2. Background of the Invention\nThe continuing use of diesel engines in motor vehicles greatly adds to the atmospheric presence of harmful pollutants such as nitrogen oxides (NOx) and particulate matter (PM). Conventional diesel engines emit NOx and/or PM substantially in excess of acceptable environmental levels. Nevertheless, because of their fuel efficiency, diesel engines remain preferable to gasoline engines for many applications. Attempts to reduce NOx and PM emissions from diesel engines have therefore continued for many years.\nFor reasons that will be discussed hereafter, the prior art has never been successful in developing a robust diesel combustion system (i.e. providing commercially acceptable responsiveness and horsepower across diverse speed and load ranges) that is capable of maintaining engine-out emissions levels of both NOx and PM simultaneously below environmentally acceptable levels. For purposes of this invention, environmentally acceptable levels of diesel NOx emissions will be defined as 0.2 g/bhp-hr or lower, and environmentally acceptable levels of diesel PM emissions will be defined as a level that can effectively be reduced to 0.01 g/bhp-hr or lower through exhaust aftertreatment.\nThe challenge of trying to maintain diesel engine emission levels of both NOx and PM simultaneously below environmentally acceptable levels has been a long-standing unsolved environmental and industry problem that has never been overcome despite extensive, concerted efforts by government and industry worldwide.\nWith the problem of simultaneously satisfactory engine-out NOx and PM reductions unsolved by the prior art for diesel engines, the diesel industry has instead turned primarily to development of NOx and PM aftertreatments (i.e. post-engine, but before the exhaust gas is released to the atmosphere) to meet current and upcoming international PM and NOx environmental regulations. However, currently many in the diesel industry view such aftertreatment systems as expensive, as retaining issues of effectiveness and durability, and as creating size (“packaging”) concerns. As an example, an ongoing challenge for NOx traps is maintaining NOx emissions always below desired levels on vehicles with diverse duty cycles, and the effectiveness and durability of NOx adsorbers may also be undermined by factors such as high temperatures or by sulfate adsorption and desulfization due to sulfur content in diesel fuel. As a second example, to effectively maintain PM levels within environmentally acceptable emission levels, PM filter traps depend on engine-out PM emissions not being too high. In addition, PM filter traps must be regenerated (i.e. cleaned, such as by combusting the particulate matter trapped therein), with the frequency of such renewal dependent upon the amount of engine-out PM the trap is forced to catch and retain.\nBecause of the above-mentioned perceived problems with the expense and performance of NOx aftertreatments in particular, it is desirable to develop an alternative, namely, a commercially acceptable diesel combustion system that can meet environmentally acceptable levels of NOx emissions (i.e. 0.2 g/bhp-hr or lower) without the need for use of NOx aftertreatments, and thereby avoid such perceived additional costs, performance, durability and packaging issues that accompany use of NOx aftertreatments.\nIn addition, because of the above-mentioned perceived problems with inconvenience and expense in regenerating PM traps, the frequency of which is dependent upon the amount by which engine-out PM emissions exceed environmentally acceptable levels (i.e. exceed 0.01 g/bhp-hr), it is also desirable to simultaneously at least maintain engine-out PM emissions low enough for PM traps to efficiently and cost-effectively be able to keep PM emissions within environmentally acceptable levels. For purposes of this application, and for ease of measuring, a smoke content level at or below a maximum of 3 BSN (Bosch Smoke Number) engine-out (i.e. before aftertreatment), preferably with an average smoke content level below 1.5 BSN, would be considered indicative of an acceptable PM engine-out exhaust level for PM aftertreatment to efficiently and repeatedly keep PM emissions within environmentally acceptable levels."}
{"text": "1. Field of the Invention\nThis invention pertains generally to the field of aircraft display units that provide information to the pilot or flight crew of an aircraft.\n2. Description of the Related Art\nHardware and software design standards for the aviation industry have been published by the Radio Technical Commission for Aeronautics (“RTCA”). The RTCA has published DO-254 entitled “Design Assurance Guidance for Airborne Electronic Hardware” and DO-178B entitled “Software Considerations in Airborne Systems and Equipment Certification.” Each publication distinguishes between simple and complex electronic hardware, recognizes five classes of failure conditions, and provides guidance for each hardware and/or software design assurance level (“DAL”) associated with a given failure condition.\nLevel A is considered the highest DAL and corresponds to a “catastrophic” failure condition of which failure of the hardware and/or software may cause a loss of the aircraft. Level B corresponds to a “hazardous” failure condition of which failure may cause a large negative impact on safety or performance, reduces the ability of the crew to operate the aircraft due to physical distress or a higher workload, or causes serious or fatal injuries among the passengers. Level C corresponds to a “major” failure condition of which failure may be significant but have a lesser impact than a hazardous failure or significantly increases crew workload. Level D corresponds to a “minor” failure condition of which failure may be noticeable but have a lesser impact than a major failure. Level E corresponds to a “no effect” failure condition of which failure has no impact on safety, aircraft operation, or crew workload.\nTo meet the higher DALs, many standards have been developed. With respect to presenting the pilot with flight-related information and aircraft-system information and/or allowing the pilot to control systems installed somewhere on the aircraft by interacting with a cockpit display unit, data exchange protocols have been published by Aeronautical Radio, Incorporated (“ARINC”) in ARINC Specification 661 entitled “Cockpit Display System Interfaces to User Systems” (“ARINC 661”), an aviation industry standard known to those skilled in the art.\nAlthough ARINC 661 could be employed for all aircraft systems, ARINC 661 may be considered too demanding for a manufacturer of an aircraft system because of the strict protocols stated therein. For example, information about real time parameter(s) of an aircraft system may be presented to the pilot and/or facilitate pilot control of the system; in order to do so, however, the strict protocols of the ARINC 661 have to be followed.\nThere may be occasions when an aircraft system does not need to meet higher DALs. For example, a cabin control system could be employed in an aircraft, but failure of such system may have no impact on safety, aircraft operation, and/or crew workload; instead, failure may mean a loss of passenger comfort. In this instance, although it may be beneficial to provide information about and/or control of the cabin control system to the pilot in the cockpit, it may not be beneficial and/or cost-effective for the manufacturer of the cabin control system to employ the strict protocols of ARINC 661 when such system does not have to meet higher DALs."}
{"text": "1. Field of the Invention\nThe present invention relates to an improvement of an axial flow fan motor for use of OA (office automation) appliance.\n2. Conventional Art\nSince most OA appliances comprise various electronic circuits within their housings. And, every electronic part consisting of those circuits generates heat. And, due to this heat generation, electronic parts of those circuits are threatened to be broken or deteriorated, so that, in this kind of OA appliances an axial flow fan motor is used for cooling. In particular, in the recent tendency of miniaturization of OA appliances, since the influence by the above heat generation becomes great, a countermeasure against the heat has become an important technical problem. And, on the other hand, the products of OA appliances such as being mass-produced have the demand for low cost from the viewpoint of easy assembly workability.\nNext, a conventional axial flow fan motor is explained based on FIG. 7, a central projection 2 is provided in a unit with a resin casing 1. Around the outer circumference of this central projection 2, iron cores 3 are fixed, around which stator windings 5 are wound through bobbins 4 made of insulating material. By this iron core 3, the bobbin 4 and the stator winding 5 the stator of the flow fan motor is formed. On the other hand, on an inner circumference of the central projection 2 bearings 6 are coupled with, and on these bearings a shaft 7 is supported. These bearings 6 are fixed by stepped portions 8 and 9 provided on the inner circumference of the central projection 2, a snap ring 10 and a spring 11. On one end of the shaft 7, a yoke 13 is fixed through a zinc casting 12 and on the inner circumference of the yoke 13, a magnet 14 is fixed. With this yoke 13 and the magnet 14, a rotor of the axial flow fan motor is formed. And on the outer circumference of the yoke 13, a resin made impeller 15 is pressed in and fixed, on which wings 16 are mounted.\nFIG. 8 shows another conventional axial flow fan motor. What is different from the conventional art shown in FIG. 7 resides in that a bearing box 17 is provided separately from the central projection 2 and adhered to the inner circumference of the central projection 2. And, by a concavo-convex engagement portion 18, the bearing box 17 is prevented from being removed. Further, this bearing box 17 is made of metal in consideration of heat resistance. The other portions of that bearing box 17 are similar with what are shown in FIG. 7, so that, by attaching the similar signs to the similar portions, explanation thereof are omitted.\nIn the above conventional art, since what is shown in FIG. 7, the central projection 2 is provided in a unitary manner in the resin made casing 1 and adapted to couple the bearings 6 with the projection 2 to support the shaft 7, the number of parts is reduced which is advantageous in assembling and invites a low cost, however, there is another problem such as mentioned below. That is, the stator consisting of the iron core 3 and the stator winding 5 may generate heat due to the copper wearing and the iron wearing. And, the iron core 3 is mounted on the outer circumference of the central projection 2 directly, so that due to the heat generated the central portion is threatened to be deformed and accelerated in its deterioration.\nFurther, by the heat too generated from the electronic parts accommodated in OA appliances the central projection is threatened to be deformed and accelerated to be deteriorated. Thus, when the central projection is deformed, the shaft center of the rotor formed with the yoke 13 and the magnet 14 is made to be displaced to cause the function of the axial flow fan motor to be deteriorated, and that, due to the deterioration of the central projection, a use life of the axial flow fan motor becomes short, by which the use life of OA appliance as a whole is threatened to become short.\nNext, in the conventional art shown in FIG. 8, since the bearing box 17 is made of a metal, there are following defects to be improved. That is, although, since the bearing box itself is made of a metal, there is no trouble against heat, due to the great thermal conductivity, the central projection 2 is heated indirectly. Accordingly, as well as in FIG. 7 which shows the conventional art, the central projection is deformed by the heat and accelerated in being deteriorated, and as a result the function of the axial flow fan motor and the use life thereof too are threatened to become short. In addition, in order to solve the problem it is recommended to make the casing 1, which is formed in a unitary manner with the central projection 2, with a resin having a heat resistance, however, at that time, a high priced resin has to be used in a great amount for that, so that the axial flow fan motor has to become high-priced.\nFurther, when the metal made bearing box 17 is inserted in the central projection 2 made of a resin, since the coefficient of thermal expansion is different between a resin and a metal with each other, there happens a case where the coupling is loosened due to the heat to be generated between different materials to be coupled. In the case, where a loosening is thus generated, since the shaft center of the rotor to the stator is resulted in being displaced, the function of the axial flow fan motor is deteriorated and a vibration is generated, which is not preferable for the OA appliance.\nStill further, since the bearing box 17 and the central projection 2 are fixed by gluing, there need a application process of a glue and a drying process thereof, which invites a difficulty of an automatic assembly to cause a decrease of productivity. In addition, since the assembly of the stator and the rotor can not be carried-out before the bearing box 17 is fixed on the central projection 2 by solidification of the glue, which causes to waste time and is not preferable from the viewpoint of assembling workability and productivity. Accordingly, one object of the present invention is to provide an axial flow fan motor, in which, in order to solve the above problems, the thermal deformation and deterioration of the bearing box are reduced to maintain the function of the axial flow fan motor, prolong the use life, reduce the use amount of a high priced resin to decrease the cost and increase the assembling workability avoiding the use of glue to increase the productivity."}
{"text": "There are many devices disclosed or proposed in the art for the application of medical materials by spraying. A particular field of spray application is for multi-part reactive materials, particularly for two-part medical materials, in which one part (“A”) is reactive with another part (“B”) either directly on contact, or upon passage of time, or upon activation, such as by heat or light. An important subfield of such applicators is dedicated to the application of fibrin sealant, in which part A contains fibrinogen, part B contains thrombin, and at least one part contains a calcium salt. On contact between the solutions, the fibrinogen is cleaved by the calcium-activated thrombin, and fibrin is formed. The fibrin self-assembles into a rigid structure, as occurs in blood clotting. This system is used to apply sealing coatings to the body, and to adhere tissue together where the mechanical stress is low. To achieve the desired effect, the two parts of the system must be intimately mixed during application to the tissue, but must be kept separate until that time. Exemplary devices for delivery of two separate fluids to a tissue surface are described in U.S. Pat. No. 5,582,596 to Fukunaga et. al., and U.S. Pat. No. 5,665,067 to Linder et. al. Linder describes a delivery device with a pair of fluid outlets surrounded by a single air annulus for spraying. Fukunaga describes a device with outlets with separate annuli, designed specifically for application of biocompatible adhesive containing human or animal protein, such as fibrin sealant, to a surgical site. The device of Fukunaga includes liquid outlets that extend distally beyond gas nozzles that surround the outlets, which Fukunaga states ensures that solutions ejected from the outlets will be sprayed uniformly by gas emitted form the nozzles.\nFibrin sealant systems can have numerous limitations, including the possible presence of viruses and prions, and newer all-synthetic systems have been devised which can alleviate these and other problems. It is easier and more reliable to tailor the mechanical properties of a synthetic system, and each of strong adherence to tissue, biocompatibility and flexibility of a coating can be achieved. An example of such a system is described in U.S. Pat. Nos. 5,749,968 and 5,800,373 and in international patent publication no. WO 96/29370. In this system, a solution of a reactive water-soluble macromer, such as a polyethylene glycol molecule extended with biodegradable groups such as lactide and made reactive by end capping with acrylate, is reacted with light in the presence of a photoinitiator or photosensitizer, and covalently polymerizes into a hydrogel. As disclosed in U.S. Pat. No. 5,748,968, it has been found that such a prior art hydrogel will adhere more strongly to tissue if the tissue is first coated with at least a part of an initiation system (herein, a “primer” or priming solution.) During polymerization, the tissue-adherent initiator assures satisfactory reaction and mechanical bonding at the tissue surface. In U.S. Pat. No. 5,800,373, it is demonstrated that this system can be further enhanced by incorporating a chemically reactive (“redox”) initiation system along with a photoinitiation system. Then one component of the redox system is used in the primer, and the other in the bulk, macromer-containing solution.\nIn these systems, and with other chemistries (e.g. as discussed in U.S. Pat. No. 5,874,500 and international patent publication no. WO 99/34833) it is often important to apply two or more fluids independently, and in some cases sequentially, to a tissue surface. Even where only one fluid is applied, it can be important to apply the fluid evenly, conveniently, and cleanly. It can also be desirable to apply the solutions in variable ratios. It can also be important to prevent any reaction of the two components in the device before the application is completed, to prevent having to clean or replace the application device during a medical procedure. It is an object of the invention to provide improved medical fluid delivery devices to address these and other needs."}
{"text": "The invention pertains to methods of bonding physical vapor deposition (PVD) target materials to backing plate materials, and also pertains to methods of bonding two aluminum-comprising masses to one another. Further, the invention pertains to structures comprising aluminum-comprising physical vapor deposition targets bonded to aluminum-comprising backing plates.\nModern developments in physical vapor deposition methodologies have created increasingly stringent requirements for robust bonding between targets and backing plates. A diagrammatic view of a portion of an exemplary sputter deposition apparatus 10 is shown in FIG. 1. Apparatus 10 comprises a backing plate 12 having a sputtering target 14 bonded thereto. A semiconductive material wafer 16 is within apparatus 10 and provided to be spaced from target 14. Sputtered material 18 is displaced from target 14 and utilized to form a coating (not shown) over wafer 16.\nAmong the modern improvements in sputter design is an increase in the distance between target 14 and semiconductive material substrate 16. Such increase in distance can enable more directional sputtering to be achieved over features of substrate 16 than can be achieved when target 14 is close to substrate 16 by allowing atoms that are not moving perpendicular to substrate 16 to land on the sidewall of the sputtering chamber. Specifically, substrate 16 will frequently have vertical holes or slots (known as vias) with depths five times their width or more (i.e., having relatively high critical dimensions). It is difficult to sputter materials into vias having high critical dimensions unless there is a relatively long throw between a sputtering target and a substrate comprising the vias.\nAlthough the longer throw creates advantages in coverage relative to shorter throw techniques, it also creates complications. One of such complications is caused by additional power utilized in long-throw technologies. The additional power can cause sputtering targets to get hotter than they had in older methods. Such heat can disrupt a bond formed between backing plate 12 and target 14. For instance, if target 14 is solder-bonded to backing plate 12, the heat developed during long-throw sputtering techniques can be sufficient to melt the solder bond and actually break target 14 free from backing plate 12. Accordingly, solder-bonding can be inappropriate for long-throw sputtering techniques.\nIt would be desirable to develop improved methods for bonding physical vapor deposition targets to backing plates.\nIn one aspect, the invention encompasses an aluminum-comprising physical vapor deposition target bonded to an aluminum-comprising backing plate to a bond strength of greater than 10,000 pounds/in2.\nIn another aspect, the invention encompasses a method of bonding a first aluminum-comprising mass to a second aluminum-comprising mass. The first aluminum-comprising mass has a first surface, and the second aluminum-comprising mass has a second surface. At least one of the first and second surfaces is exposed to a mixture comprising HF and HNO3, and the first surface is then contacted with the second surface. The first and second masses are pressed together to bond the first mass to the second mass.\nIn yet another aspect, the invention encompasses a method of bonding a physical vapor deposition target material to a backing plate material. The target material and backing plate material are joined in physical contact with one another. The target material and backing plate material are then compressed under a load that progresses sequentially as follows: (1) the load is maintained at a first pressure of at least about 4 ksi for a time of at least about 20 seconds; (2) the load is increased by at least about 0.5 ksi to a new pressure; and (3) the load is maintained at the new pressure for at least about 20 seconds."}
{"text": "1. Field of the Invention\nThe present invention relates to thermostats. More specifically, the present invention is directed to a plural zone control system using respective zone thermostats.\n2. Description of the Prior Art\nConventional heating and cooling systems for residential buildings are typically controlled with a single thermostat. Thus, only one temperature can be specified by the homeowner. The temperature in the vicinity of the thermostat will be controlled to the desired level, but in other parts of the house the temperature can vary widely due to inadequate air distribution, solar radiation entering through south facing windows, outdoor wind, heat generated by people and appliances and many other factors. In response to these problems and the desire for greater comfort and flexibility, so-called \"zoned\" systems have been developed. A zoned system allows different parts of a residence to be controlled at different temperatures or at equal temperatures. This is accomplished by a separate thermostat in each temperature zone and a means for controlling the flow of heat or cool to that zone. The most common heating/cooling system is the so-called forced warm air system. In a zoned system of this type, dampers are placed in the ductwork to allow the airflow to each zone to be turned on or off or to be modulated.\nAlthough the zoned heating/cooling system described above offers the advantages of greater comfort through the more accurate control of temperature, greater flexibility by allowing the homeowner to set different temperatures through the house and the potential for energy savings by discontinuing the heating or cooling of unoccupied areas of the home, these systems also have some disadvantages. In its most common form, a zoned system allows each individual thermostat to turn on the heating/cooling plant (furnace, air conditioner, etc.) and simultaneously open the corresponding zone damper while closing the dampers to other zones not requiring heating/cooling. While this approach is quite simple and easy to implement in practice, it creates the following problems. With several individual thermostats independently turning the heating/cooling plant on and off, this appliance will experience significantly more cycles than normal and the service life of the plant will therefore be shortened. For example, without coordinating the calls for heating or cooling, the potential exists for short-cycling the plant, e.g., thermostat A might shut off the furnace, but thermostat B turns it back on 15 seconds later). These short cycles can damage the heating/cooling plant.\nFurther, when only one zone thermostat requests or calls for heat, the entire output of the heating plant (which is normally capable of heating the entire house) will be supplied to that zone causing large fluctuations in air temperature. When only one zone thermostat calls for heat, the air flow rate through the heating plant may be reduced because the ductwork to that zone cannot conduct the full flow capacity of the blower. This may lead to excessive temperatures in the plant and may cause the high temperature safety switch to trip repeatedly. On the other hand, during cooling operation, the reduced airflow can lead to ice formation on the cooling heat exchanger and similar safety and reliability problems. The present invention is directed to a zone control system that addresses the problems discussed above by allowing multiple thermostat requests to be coordinated such that the air flow through the system is maximized, and the on/off cycles of the plant are minimized."}
{"text": "The present invention relates generally to a bicycle saddle, and more particularly to a bicycle saddle supporting frame which is provided with a bracing member for preventing the bicycle saddle supporting frame from moving sideway at the time when the saddle is loaded with a weight.\nAs shown in FIG. 1, the conventional bicycle saddle has a main body 1, which has a profile becoming progressively wider from the front end thereof toward the rear end thereof and is supported by a metal frame 2 in such a manner that the underside of the front end of the main body 1 is fastened with a curved portion 3 of the metal frame 2, and that the underside of the rear end of the main body 1 is fastened with two gradually-diverging arms 4. The main body 1 is mounted on a seat post 6 in conjunction with a connection seat 5 which is fastened at one end thereof with the two arms 4 and at other end thereof with the top end of the seat post 6.\nSuch a conventional form of fastening the bicycle saddle with the bicycle seat post 6 as described above is defective in design in that the metal frame 2 has a propensity to sway to one side, as shown by the dotted lines in FIG. 2, at such time when the corresponding side of the metal frame 2 is exerted on by an external force as shown by an arrow in FIG. 2. It is conceivable that the lateral movement of the bicycle saddle can cause a great deal of discomfort to the hips of a bicyclist. In addition, the lateral movement of the metal frame 2 can result in damage or even crack of fastening portions 8 between the main body 1 and the two arms 4 of the metal frame 2.\nThe primary objective of the present invention is therefore to provide a simple and cost-effective technique for overcoming the deficiencies of the conventional bicycle saddle as described above.\nThe technique of the present invention involves a measure by which two gradually-diverging arms of the bicycle saddle supporting frame are reinforced by a bracing means, which is preferably made of a rigid material. The two arms are reinforced by the bracing means such that the two arms are slightly expandable in opposite directions so as to be fastened with both sides of the underside of the main body of the bicycle saddle. In addition, the bracing means of the present invention is compatible with the bicycle saddles of various sizes in view of the fact that the bracing means enables the distance between the two retaining portions of the two arms to be adjusted.\nThe objective, features, functions, and advantages of the present invention will be more readily understood upon a thoughtful deliberation of the following detailed description of two preferred embodiments of the present intention with reference to the accompanying drawings."}
{"text": "Current methods for increasing the likelihood of successful incremental re-training of a predictive model exhibit a plurality of problems that make current systems insufficient, ineffective and/or the like. Through applied effort, ingenuity, and innovation, solutions to improve such methods have been realized and are described in connection with embodiments of the present invention."}
{"text": "The present invention additionally provides multistage coating methods using these coating material compositions, and also the use of the coating material compositions as clearcoat material, and application of the coating method for automotive refinish and/or for the coating of plastics substrates and/or of utility vehicles.\nPolyurethane coating materials typically comprise a catalyst, and in this context not only acidic compounds but also, in particular, tertiary amines and/or metallic compounds, such as various tin compounds, for example, more particularly dibutyltin dilaurate and dibutyltin oxide, are employed.\nThe employment of tin-containing catalysts is to be avoided in coating materials, as elsewhere, because of the toxicity inherent in many tin compounds. The EU Commission's Working Group on Classification and Labelling” have categorized dibutyltin oxide (DBTO) and dibutyltin dilaurate (DBTL) accordingly.\nThe article available on the Internet at the address www.wernerblank.com and titled “Catalysis of the Isocyanate-Hydroxyl Reaction by Non-Tin Catalysts”, by Werner J. Blank, Z. A. He, and Ed. T. Hessell of the company King Industries Inc., therefore describes alternatives to the typical tin-containing catalysts based on various metal salts and metal complexes, such as zirconium chelates, aluminum chelate, and bismuth carboxylate.\nDE 10 2008 061 329 A1 discloses coating materials where the use of metal-containing catalysts is to be avoided as far as possible and which instead as catalyst comprise 1,3 substituted imidazolium salts for the deblocking of blocked polyisocyanates in polyurethane coating materials.\nWO04/029121 describes polyurethane compositions which are stabilized in terms of the reactivity of the composition by addition of acids with a pKa range between 2.8 and 4.5, these acids being able to be utilized at the same time as catalyst. Acids used in this context and with a pKa range between 2.8 and 4.5 include, for example, benzoic acid, hydroxybenzoic acid, salicylic acid, phthalic acid, and so on. The compositions preferably comprise no further catalyst, although in addition it is also possible to use the typical known polyurethane catalysts, such as tertiary amines or amidines or organometallic compounds, such as tin compounds more particularly. Where amines are used as catalyst, it is necessary to employ great care in the selection of the type of amine and its amount, since the aminic catalysts are able in part to eliminate the stabilizing action of the organic acids added.\nU.S. Pat. No. 5,847,044 describes polyurethane powder coating materials which as catalysts comprise N,N,N′-trisubstituted amidines, more particularly bicyclic amidines.\nWO 09/135600 describes polyurethane compositions, more particularly sealants, adhesives, and foams, which comprise as catalyst the reaction product of a metal salt with nitrogen-containing, heterocyclic compounds, more particularly substituted imidazoles.\nDE-A-24 34 185 describes a process for preparing amidine-metal complexes and their use as catalysts for the isocyanate polyaddition reaction. These amidine-metal complexes are prepared by reacting an amidine with a 0.5- to 4-fold molar amount of a metal compound, the amidines used comprising not only monocyclic and/or bicyclic compounds, such as imidazoles more particularly, but also acyclic compounds, such as formamidines, acetamidines, benzamidines, and guanidines. Metal compounds used are those of trivalent iron, of divalent nickel, of divalent zinc, of divalent manganese, of divalent tin or of tetravalent tin, with the corresponding carboxylates being employed more particularly.\nLastly, U.S. Pat. No. 7,485,729 B2 and also the equivalent specifications WO06/022899, US 2006/0247341 A1, and US 2009/0011124 A1, describe organometallic compounds and coating materials comprising them. Coating materials described are powder coating materials based on hydroxyl-containing polyacrylates and/or polyesters and on uretdione-containing polyisocyanates, liquid coating materials based on hydroxyl-containing polyacrylates and/or polyesters and on blocked polyisocyanates, and also solventborne coating materials based on epoxy/carboxy or epoxy/anhydride components. The organometallic compounds used as catalyst, besides other metal-amidine complexes, are cyclic or acyclic zinc biscarboxylate-bisamidine complexes, such as Zn(1,1,3,3-tetramethylguanidine)2(2-ethylhexanoate)2, for example.\nA problem addressed by the present invention, therefore, was that of providing coating material compositions, more particularly for automotive refinish and for the coating of utility vehicles, that ensure good assembly strength after just a very short time, meaning that they ought to ensure rapid curing even under the conditions of refinish and of the finishing of utility vehicles, in other words ought after curing at 60° C. for 30 minutes already to have undergone curing to an extent such that initial assembly operations or demasking operations can be carried out without damage to the coating. At the same time, however, the coating material compositions ought, at room temperature and after mixing of the binder component with the isocyanate component, to have a good potlife of at least 2 hours. Potlife here means the period of time within which the coating material composition has attained twice its initial viscosity. Moreover, the coating material compositions ought to lead to coatings exhibiting good through-curing and sufficient ultimate hardness. Furthermore, these coating material compositions ought not to show any changes in color before and after curing. Particularly in the field of clearcoat materials in the automotive industry, the intrinsic color of the systems is subject to cracking requirements. Thus the catalyst neither must exhibit any intrinsic color and nor must it lead to discoloring at mixing or during curing of the coating material when the catalyst is mixed with the typical components of a coating material. The resultant cured coatings ought, furthermore, to have no tendency toward yellowing after exposure in the test known as the WOM test (WOM=Weather-Ometer Test, determined in accordance with SAE (Society of Automotive Engineers) Standard J2527_04). Yellowing is determined using the multiple angle colorimeter BYK-mac from BYK-Gardner GmbH, D-82538 Geretsried, with calculation according to DIN 6174.\nFurthermore, the catalyst ought to be able to be added to the coating system from the outset. However, admixing the catalyst to the coating systems from the outset is not to cause any adverse effect on the shelflife of the coating composition. Furthermore, the catalyst ought to be insensitive to hydrolysis, since even in systems in organic solution, the typically high concentration of hydroxyl groups can result in a reduction in catalyst activity over the storage period. Especially in the automotive refinish segment, an extremely long shelflife even at relatively high temperatures is an advantage.\nLastly, the coating material compositions ought to be able to be prepared simply and with very good reproducibility, and ought not to cause any environmental problems during application. More particularly, catalysts containing tin ought to be avoided or at best be entirely dispensible."}
{"text": "From U.S. Pat. No. 6,789,106, it is known a method and a structure that allows a user who has access to an audio/visual object to cause an interactive control element of an interactive multimedia system to grab and store for future retrieval and peruse the audio/visual object by simply activating a selection mechanism of the interactive multimedia system. The interactive control element of the interactive multimedia system receives a plurality of audio/visual objects for play from the content provider and causes each such object to be stored to a temporary storage element as it is played. In response to the user activating the selection mechanism of the interactive multimedia system, it is determined whether the interactive control element has completed storing an audio/visual object to the temporary storage element that the user has selected via the selection mechanism. If the interactive control element has completed storing the selected audio/visual object to the temporary storage element, the selected audio/visual object is transferred from the temporary storage element to a long-term storage element. If, however, the interactive control element has not completed storing the selected audio/visual object to the temporary storage element, it is determined whether the selected audio/visual object is currently being played. If it is not currently being played, this means that the selected object may not be available locally and the interactive control element retrieves the selected audio/visual object from the content provider (if necessary) and stores it in the long-term storage element. If the selected audio/visual object is being currently played, then the interactive control element continues to store the selected audio/visual object within the temporary storage element until the selected audio/visual object is stored in its entirety (has finished playing) and then transfers it from the temporary storage element to the long-term storage element, where it can be retrieved at some later time when desired.\nFrom patent application no US2002/0087988, it is known an interactive TV contextual content surfing method, system and program product that enables viewers to surf from one related program to another program through hyperlinking. A TV broadcast station receives interactive TV program content generated by creators. The content is categorized by the creator and provided to the broadcast station in the MPEG or webstreaming video/audio formats. A broadcast server stores a dynamic table of correspondence between TV program categories and TV channel numbers. The hyperlinking information contains a list of channels that broadcast the same nature of programs by referencing the correspondence table. The server broadcasts the hyperlinking information to the list of channels. The broadcast network carries the TV programs and associated classification information in separate channels to the local TV station for retransmission to TV devices including set-top boxes. The set-top box stores program information by categories in classification tables, contents of a TV program guide section, and recorded programs.\nFrom U.S. Pat. No. 6,215,483, it is known a method and apparatus for receiving logical address links in advance of a television program and combining such links with those links that may be received in real-time during a television broadcast. According to one aspect of this invention, content associated with a TV program may be automatically displayed by a client system, such as an Internet terminal or a computer equipped with a TV tuner. Prior to broadcast of a TV program, TV listing information is stored in the client system. The TV listing information includes one or more batch mode logical addresses, e.g. Uniform Resource Locators (URLs), that identify content, such as Internet content. The client system alerts the viewer of the existence of an active logical address associated with the current TV program by providing a visual indication. After the visual indication has been selected by the viewer, the client system displays the content identified by the logical address. Additionally, one or more real-time logical addresses may be received by the client system embedded in a video signal associated with the TV program. The real-time logical addresses may be embedded in either a text service channel (e.g., T1, T2, T3, T4) or a captioning service channel (e.g., CC1, CC2, CC3, CC4) of the vertical blanking interval (VBI) of the video signal. When both real-time logical addresses and batch mode logical addresses are associated with a TV program, the client system determines which to bring to the attention of the user of the client system based upon a set of predefined rules. Features of the present invention are applicable to many broadcast television (TV) systems including National Television Standards Committee (NTSC), Phase Alternate Lines (PAL), and Sequential Couleur Avec Memoire (SECAM) as well as the proposed High Definition Television (HDTV) standard.\nFrom patent application no WO02/085024, it is known a video-processing method and system for enriching one or several primary video signals by association of additional data in a video scene encoded in accordance with the MPEG4 standard, this associating allowing simultaneous display of a primary video signal and said additional data, and enabling a user to interact with and access the contents of said additional data.\nFrom patent application no US2005/0193425, it is known a method for delivery and presentation of information relevant to the contents of frames in an audio/video program that enables TV viewers to retrieve information on the contents (for example, objects, items, concepts and the like) contained in a frame or a set of frames (video segments) when they watch TV or video programs. The information relevant to the contents of frame(s) is delivered to a STB or DVR by third-party service providers through back channels such as the Internet if the information of how to accurately access the frames pointed by STB users are delivered to the service providers, and the content-relevant information may be presented in the form of a GUI for the TV viewer. This disclosure relates to the processing of program guide information and, more particularly, to techniques for delivering information on video segments of broadcast TV programs to STBs."}
{"text": "There have usually been used a plastic microplate provided thereon with a plurality of wells as a reaction vessel or a culture vessel in a chemical test which comprises a process for reacting a chemical sample with a certain reagent, or a microbiological test which comprises a step for cultivating a microorganism in a sample, in particular, in a test which requires treatments of a number of samples, reactions with a number of reagents or cultivation of a microorganism in a number of systems.\nFor instance, a plastic plate equipped with a number of small wells has usually been used in immunological assays such as ELISA assays. In respect of the determination of the sensitivity of bacteria to antibacterial agents by the micro liquid dilution technique, i.e., the determination of the minimum inhibitory concentration (MIC) for a specific antibacterial agent, there has widely been used the single disc method according to the agar plate dilution method or the agar diffusion method, but the use of a microplate provided with U-shaped wells identical to that used in the immunological tests is specified, in the standard method for the determination of MIC by the micro liquid dilution technique which is designated by the Chemotherapeutic Association of Japan in 1989, which is similar to the foregoing single disc method further economized and automated.\nIn such a test wherein a microplate is used, individual reactions or cultivation of microorganisms each is carried out in a separate well on the plate. In the determination of MIC by the micro liquid dilution method using a microplate, for instance, an antibacterial agent is added to a plurality of wells in a variety of concentrations and the growth of a bacterium to be tested in each well is then observed to thus determine the minimum inhibitory concentration of the antibacterial agent. As microplates used in such tests, there have been put on the market one which is prepared by preliminarily dispensing a desired amount of an antibacterial agent to each well and then drying or freezing and storing. However, the microplate per se is expensive, this makes it difficult to offer a cheaper commodity and accordingly, the users have strongly desired for the reduction of the price thereof. Moreover, such a microplate requires the dispensation of a desired amount of, for instance, a liquid sample to each well during testing and it requires a tremendous labor to carry out the test. Accordingly, it is preferred to develop an instrument for easily dispensing, for instance, a desired amount of a liquid sample to a container used for carrying out the desired test.\nTo meet such a demand, the inventors of this invention have proposed an instrument for use in chemical tests or microbiological tests, which comprises a base provided thereon with portions for holding a sample (WO 94/04703). The instrument is designed to absorb, for instance, a liquid sample in the sample-holding portion formed on the base to hold a desired amount of the liquid sample on or within the sample-holding portion and to thus meet the foregoing requirement. To allow the instrument to hold a sample, however, it should be designed in such a manner that an excess of the sample is dropwise added to the instrument and then it is, for instance, inclined to remove the excess sample; or the instrument per se is immersed in a liquid sample accommodated in a container and then pulled up to thus allow the sample-holding portion to hold a desired amount of the sample.\nFor this reason, the liquid sample is inevitably scattered in the surrounding to some extent when using the foregoing instrument to carry out such a test. However, any scattering of the liquid sample should sometimes be avoided when the liquid sample to be tested comprises, for instance, pathogens."}
{"text": "Electrical cables for transmission of electrical signals are known. One common type of electrical cable is a coaxial cable. Coaxial cables generally include an electrically conductive wire surrounded by an insulator. The wire and insulator are surrounded by a shield, and the wire, insulator, and shield are surrounded by a jacket. Another common type of electrical cable is a shielded electrical cable comprising one or more insulated signal conductors surrounded by a shielding layer formed, for example, by a metal foil. To facilitate electrical connection of the shielding layer, a further un-insulated conductor is sometimes provided between the shielding layer and the insulation of the signal conductor or conductors. Both these common types of electrical cable normally require the use of specifically designed connectors for termination and are often not suitable for the use of mass-termination techniques, i.e., the simultaneous connection of a plurality of conductors to individual contact elements, such as, e.g., electrical contacts of an electrical connector or contact elements on a printed circuit board. Although electrical cables have been developed to facilitate these mass-termination techniques, these cables often have limitations in the ability to mass-produce them, in the ability to prepare their termination ends, in their flexibility, and in their electrical performance. In view of the advancements in high speed electrical and electronic components, a continuing need exists for electrical cables that are capable of transmitting high speed signals, facilitate mass-termination techniques, are cost-effective, and can be used in a large number of applications."}
{"text": "1. Technical Field\nThe present invention relates to a method and system for data processing in general, and in particular to a method and system for providing automatic notification. Still more particularly, the present invention relates to a method and system for providing automatic notification of an end of lease of a computer and its location.\n2. Description of the Prior Art\nMore and more companies prefer computer leasing over purchasing. The lease term of a typical computer lease usually lasts approximately two, three, or four years. At the end of a computer lease, a company needs to return all the lease computers; otherwise, the company must continue to pay at the lease rate for any lease computers that have not been returned. The task of finding all the lease computers is usually quite straightforward when there are only a few lease computers. However, when there is a large number of lease computers, such as several thousand units, installed at various locations within a company, then the task of finding all the lease computers becomes daunting.\nA database can be used to keep track of information such as whether or not a computer lease for a specific computer is about to expire, along with the location of the computer according to its first installation. However, computers tend to be relocated many times within a huge organization after their first installation, and yet the database is usually not updated after the relocation of computers. As a result, the database tracking method becomes very impractical. Consequently, it would be desirable to provide an improved method for providing automatic notification when a computer lease is about to expire, and for providing the location of the computer to which the computer lease belongs."}
{"text": "Resistance to infection by microorganisms involves both non-specific mechanisms and adaptive immune responses mediated by B and T lymphocytes. Non-specific mechanisms such as enzyme action, pH and epithelial tissue secretions prevent invasion by many organisms. However, when the first line of defense is breached, phagocytes engulf infectious agents and process them leading to the stimulation of B and T cells. Deficiencies in the phagocytic system, whether hereditary or acquired, have serious consequences, since microorganisms with low pathogenicity for normal individuals can cause recurrent infections in individuals with ineffective phagocytic systems.\nPolymorphonuclear leukocytes (PMNs) are of particular importance for combating infections. These cells contain an enzyme, myeloperoxidase, with well documented microbicidal action. PMNs act non-specifically by phagocytosis to engulf microorganisms, incorporate them into vacuoles, termed phagosomes, which fuse with granules containing myeloperoxidase to form phagolysosomes. In phagolysosomes the enzymatic activity of the myeloperoxidase leads to the formation (from hydrogen peroxide and chloride) of hypochlorous acid, a potent bactericidal compound. Hypochlorous acid is oxidizing in itself, and reacts most avidly with thiols and thioethers but also converts amines into chloramines and chlorinates aromatic amino acids. Macrophages are large phagocytic cells which, like PMNs are capable of phagocytosing microorganisms. Macrophages can generate hydrogen peroxide and upon activation also produce myeloperoxidase. In addition, myeloperoxidase in the plasma is taken up by macrophages.\nLinkage of myeloperoxidase activity to disease has been implicated in numerous inflammatory diseases including Alzheimer's disease, multiple sclerosis, asthma, atherosclerosis, cancer, cystic fibrosis, chronic obstructive pulmonary disease, inflammatory bowel disease, and rheumatoid arthritis. Assays for identification of inhibitors of myeloperoxidase will be invaluable in demonstrating the contribution of hypochlorous acid to the pathogenesis of such diseases, as well as for diagnostic tools. They will also provide the platform for the development of effective drugs against myeloperoxidase-dependent tissue damage. A suitable screening assay for MPO inhibitors should measure the chlorination activity of myeloperoxidase under realistic physiological conditions. It must be simple, sensitive, and produce a coloured or fluorescent product, and it must not perturb the activity of MPO. The detector will need to be present at high enough concentrations to scavenge all the hypochlorous acid produced but not so as to react directly with the enzyme. Currently, the most commonly used assay is based on the primary formation of taurine chloramine from hypochlorous acid and taurine. Previous attempts to use modified forms of this taurine chloramine assay for measuring the activity of myeloperoxidase have revealed specific problems that negate its use for screening potential inhibitors of the enzyme. Major problems are that dimethylsulphoxide (DMSO), the solvent in which test compounds are most conveniently and usually dissolved, interferes with the assay, and many of the test compounds absorb in the same region as the commonly used chromophore.\nPhysiologically, PMN-derived myeloperoxidase primarily catalyzes the hydrogen peroxide-dependent oxidation of halides such as chloride ions, but its ability to use electron donors such as O-dianisidine as a substrate forms the basis of an assay to quantify PMNs in inflamed tissue. Adaptations of the method to avoid artifacts caused by tissue/blood constituents, such as the cytosolic enzymes, catalase and glutathione peroxidase, reducing substrates such as glutathione and ascorbic acid, and heme proteins have been described (Grisham et al., Methods in Enzymology, Vol. 186, Ed. Packer et al., San Diego: Academic Press, pp. 729–742 (1990)). However, there is a need for a more reliable analytical method of quantifying the PMN content of inflamed tissues that is not affected by other endogenous tissue constituents.\nThe present invention describes a simple, sensitive and robust assay that avoids the major problems associated with currently known assays and that is well suited to the identification of inhibitors of myeloperoxidase. In another aspect, the present invention is useful for certain diagnostic tests associated with myeloperoxidase."}
{"text": "1. Field of the Invention\nThe present invention concerns an apparatus to control an antenna arrangement in a magnetic resonance apparatus.\n2. Description of the Prior Art\nMagnetic resonance apparatuses with a transmission antenna for multi-channel measurements are increasingly gaining importance.\nFIG. 1 shows a first activation of an antenna arrangement ANT1 according to the prior art.\nThe antenna arrangement ANT1 here is fashioned as a cylindrical birdcage antenna and serves for acquisition of a patient in a magnetic resonance examination.\nThe antenna arrangement ANT1 has eight longitudinal rods LS11 through LS18. Both ends of the longitudinal rods LS11 through LS18 are connected with one another by respective circular termination rings AR11 and AR12. Each termination ring has capacitors in the circumferential direction as needed.\nThe shown antenna arrangement ANT1 is designed for a multi-channel system so that each longitudinal rod LS11 through LS18 has an associated infeed point AP11 through AP18. Amplified transmission signals for emission are respectively connected at each of the infeed points AP11 through AP18. The transmission signals differ from one another.\nOnly one infeed situation at a first infeed point AP11 of a first longitudinal rod LS11 is described in the following as an example.\nThe infeed situation at the infeed points AP12 through AP18 is correspondingly designed.\nThe first longitudinal rod LS11 serves as an antenna element and has the first infeed point AP11 with two connections AS11 and AS12. An amplified transmission signal SV11 for emission is connected at the two connections AS11 and AS12.\nThe amplified transmission signal SV11 is formed by an amplifier RFPA1 at whose input a radio-frequency signal SS11 arrives via a signal line SL11.\nThe amplifier RFPA1 is connected at the output side via two feed lines ZL11 and ZL12 with the two connections AS11 and AS12 of the infeed point AP11.\nAs an alternative (not shown), it is possible for a number of the infeed points to be arranged on one of the termination rings at which transmission signals of different amplification are connected for a multi-channel system for emission.\nThe type of antenna excitation shown here is designated as a “single rod excitation”.\nFIG. 2 shows a second control of an antenna arrangement ANT2 according to the prior art.\nThe antenna arrangement ANT2 here is fashioned in an arc and has a curved antenna element SE21 with an infeed point AP21.\nThe infeed point AP21 has two connections AS21 and AS22 at which an amplified transmission signal SV21 is connected for emission.\nThe amplified transmission signal SV21 is formed by an amplifier RFPA2 at whose input a radio-frequency signal SS21 arrives via a signal line SL21.\nThe amplifier RFPA2 is connected at the output side via two feed lines ZL21 and ZL22 with the two connections AS21 and AS22 of the infeed point AP21.\nHowever, the two embodiments described have disadvantages, as follows.\nThe amplified transmission signals are presently directed from the amplifier outputs to the infeed points via relatively thick, low-attenuation coaxial lines of specific capacity. Due to the unwieldy coaxial cables, this solution is difficult to manage and is also error-prone, depending on the wiring arrangement.\nThe transmission signals to be amplified are correspondingly also directed via coaxial lines to the associated amplifiers.\nIt should be noted that additional signal attenuations are caused even given use of low-attenuation coaxial cables."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a production method for a metal matrix composite material. More specifically, the present invention relates to a production method for a metal matrix composite material excellent in properties, such as plastic workability, thermal conductivity, room-temperature or high-temperature strength, high stiffness, neutron absorption performance, wear resistance and low thermal expansibility.\n2. Description of the Related Art\nHeretofore, there has been known a method of producing a composite material having an aluminum matrix through a powder metallurgy process, comprising the steps of:\n(1) mixing a powder of a ceramic material serving as a reinforcing material, such as Al2O3, SiC, B4C, BN, aluminum nitride or silicon nitride, with an aluminum powder serving as a matrix;\n(2) subjecting the mixed powder to canning or cold compaction to form a compact;\n(3) subjecting the compact to degassing, sintering, etc.; and\n(4) forming the sintered compact into a desired shape.\nThe sintering process in the step (3) includes: a technique (A) of simply heating the compact; a technique (B) of pressing the compact at high temperatures, such as hot pressing; a technique (C) of sintering the compact through hot plastic working, such as hot extruding, hot forging or hot rolling; a technique (D) of pressing the compact while applying a pulse current thereto, i.e., subjecting the compact to so-called “pulse-current pressure sintering” (as disclosed, for example, in JP 2001-329302A); and a technique (E) based on a combination of two or more of the techniques (A) to (D). There has also been known a technique of performing the sintering process in conjunction with the degassing process.\nIn recent years, aluminum matrix composite materials have been increasingly developed for use in new applications requiring not only strength but also a high Young's modulus, wear resistance, low thermal expansibility and nuclear-radiation absorption capability. Although each of the functions can be enhanced by increasing an amount of ceramic additives having the respective functions, an approach of simply increasing the amount of ceramic additives will cause significant deterioration in sinterability and plastic workability, such as, extrudability, rollability or forgeability.\nFrom this standpoint, there has been proposed a technique of preparing a ceramic preform, and impregnating the ceramic preform with molten aluminum alloy to allow ceramic particles to be uniformly dispersed over an aluminum alloy matrix in a high density. In reality, this technique is likely to involve problems about insufficiency of the impregnation with the molten aluminum alloy, and occurrence of defects, such as shrinkage during solidification of the molten aluminum alloy.\nInternational Publication No. WO 2006/070879 relates to a method of producing an aluminum matrix composite material, which comprises the steps of: (a) mixing an aluminum powder and a ceramic powder to prepare a mixed powder; (b) subjecting the mixed powder to pulse-current pressure sintering together with a metal sheet to form a cladded material where a sintered compact is cladded with the metal sheet; and (c) subjecting the cladded material to plastic working to obtain an aluminum matrix composite material.\nIn WO 2006/070879, before a mixed powder prepared by mixing an aluminum powder and a ceramic powder is subjected to a rolling process, it is necessary to subject the mixed powder to pulse-current pressure sintering while sandwiching the mixed powder between metal sheets, so as to preform a cladded material in such a manner as to be maintained in a predetermined shape. The reason is that it is difficult or substantially impossible to roll the cladded material unless it is performed by sintering in such a manner as to be maintained in a predetermined shape.\nAs described above, it is essential for WO 2006/070879 to preform the cladded material in such a manner as to be maintained in a predetermined shape, i.e., to subject the mixed powder to pulse-current pressure sintering, which leads to deterioration in process efficiency and difficulty in achieving an intended cost reduction. Thus, there remains a strong need for solving these problems."}
{"text": "1. Field of the Invention\nThe present invention relates to a process and a circuit for regulating welding current and power as a function of welding speed in welding devices for continuous longitudinal-seam welding of metal bands formed into a tube.\n2. Description of the Prior Art\nA circuit arrangement is already known (German OS 19 00 856 corresponding to U.S. Pat. No. 3,619,554 ) in which the electrodes of a welding device are affected from a voltage source via a transformer, a control element adjusting the welding current, and a speed-dependent tachogenerator. In this case, means is provided for controlling the magnitude of the welding current as a function of the supply voltage, the electrode distance and temperature variations.\nIn a control circuit of this type, for DC welding, the actual welding current is continuously measured and, by means of disturbance-variable compensation, the control element is affected in the case of even very small variations, so that the welding current corresponds to the requirements in each case. The adjustment in this case can be carried out promptly, so that variations in the supply voltage, the electrode distance or temperature changes are not noticeable with regard to the welding current, but the magnitude of the welding current always takes on the value required for thorough welding of the weld in each case.\nHowever, problems arise when higher manufacturing speeds are required and different metals or band thicknesses are used. This is because high manufacturing speeds permit the disturbing effect of the waviness of the welding current to increase, so that it is no longer assured that satisfactory welds can be produced with known control devices.\nTo provide a remedy in this case, it would be possible to incorporate filter elements in the control circuit, but there are limits to a significant enlargement of such filter networks because of the time constants of the control circuit. The replacement of the known control elements, for example by transistors, also does not always produce the desired effect or makes additional expenses necessary, for example, by the fact that the excess energy must be removed by means of suitable coolants, such as water or air."}
{"text": "1. Field of the Invention\nThe present invention relates to a highly reliable solar cell module. More particularly, the present invention relates to a highly reliable solar cell module comprising a photovoltaic element (a solar cell) as a photoelectric conversion member and a specific surface side covering material disposed to seal a front side (a light receiving face) of said photovoltaic element, which excels particularly in weatherability and moisture resistance and has good exterior appearance, and can be used as a building construction material such as a roofing material.\nThe present invention includes a process for the production of said solar cell module.\n2. Related Background Art\nThere have been proposed a number of solar cell modules comprising a photovoltaic element (a solar cell) sealed by a resin, a surface protective member which covers a front side (a light receiving side) of said photovoltaic element, and a back side protective member which covers a back side of said photovoltaic element.\nFIG. 5 is a schematic cross-sectional view illustrating the constitution of an example of such solar cell module.\nIn FIG. 5, reference numeral 501 indicates a transparent surface protective member, reference numeral 502 a transparent thermoplastic resin as a filler, reference numeral 503 a photovoltaic element (or a solar cell), reference numeral 504 an insulating back face protective member. In the solar cell module shown in FIG. 5, the photovoltaic element 503 is enclosed by the transparent thermoplastic resin 502, the front side (the light receiving side) of the photovoltaic element 503 is covered by the transparent surface protective member 501, and the back side of the photovoltaic element 503 is covered by the insulating back face protective member 504.\nThe surface protective member 501 comprises a transparent resin film made of a fluororesin such as ETFE (ethylene-tetrafluoroethylene copolymer) or PVF (polyvinyl fluoride), or a glass plate. The thermoplastic resin 502 comprises EVA (ethylene-vinyl acetate copolymer) or butyral resin. The back face protective member 504 comprises an organic film such as a nylon film or a composite comprising an aluminum foil sandwiched with TEDLAR (trademark name).\nThe back face protective member 504 is used in order to reinforce the solar cell module while adding an appropriate rigidity thereto.\nThe thermoplastic resin 502 serves as a filler enclosing the photovoltaic element 503 so as to prevent the photovoltaic element from being externally damaged and from external shock. The thermoplastic resin 502 also serves not only as an adhesive between the photovoltaic element 503 and the surface protective member 501 (that is, the fluororesin film) but also as an adhesive between the photovoltaic element 503 and the back face protective member 504.\nNow, the filler (comprising the thermoplastic resin) is used for covering irregularities of the photovoltaic element while ensuring the adhesion of the photovoltaic element with not only the surface protective member but also the back face protective member. Therefore, it is required for the filler to have sufficient weatherability, adhesion properties, and heat resistance. Besides, the filler is also required to excel in long-term durability, endure against thermal expansion and heat contraction, and have flexibility.\nAs the filler, any transparent thermoplastic resins can be optionally used as long as they satisfy the above requirements. However, it is advantageous to use ethylene series copolymers because they excel particularly in heat resistance and adhesion properties and are relatively inexpensive. Of these ethylene series copolymers, EVA (ethylene-vinyl acetate copolymer) is the most often used.\nIn the case where the solar cell module is used under an environment with high temperatures, in order to ensure the adhesion of the filler (comprising the thermoplastic resin) with the photovoltaic element, the surface protective member, and the back face protective member under such severe environment, the filler is usually crosslinked with a crosslinking agent. Particularly in the case of using EVA as the filler, the EVA is usually crosslinked with a crosslinking agent comprising an organic peroxide.\nHowever, in this case, such problems as will be described are liable to occur. The organic peroxide used as the crosslinking agent is liable to decompose an antioxidant added to the thermoplastic resin as the filler, resulting in deteriorating the weatherability of the filler. This situation becomes a serious drawback in the case where the thermoplastic resin as the filler situated on the light incident side (the front side) of the photovoltaic element is crosslinked with the organic peroxide as the crosslinking agent.\nBesides, in the case where a solar cell module having a filler comprising a given thermoplastic resin such as EVA containing an organic peroxide as a crosslinking agent is prepared by means of vacuum lamination, sometimes the organic peroxide is not sufficiently removed and remains in the form of a bubble in the filler, whereby markedly deteriorating the bonding property and moisture resistance of the filler. In addition, gas is generated when the filler is crosslinked with the organic peroxide, and this gas tends to form a bubble in the filler, resulting in deteriorating the bonding property and moisture resistance of the filler. The formation of this bubble is more likely to occur upon rapid temperature rise.\nFurther, a relatively long period of time and a relatively large quantity of energy are required in order to crosslink the thermoplastic resin as the filler. This situation is a factor of making the resulting solar cell module costly.\nFor a photovoltaic element comprising a semiconductor layer as a photoelectric conversion member formed on an electrically conductive metal substrate and a transparent electrically conductive layer formed on said semiconductor layer which is used in a solar cell module, it is liable to have a short-circuit defect (a shunt) in which the substrate is short-circuited with the transparent electrically conductive layer due to irregularities present at the surface of the substrate or/and unevenness in the film formation. This short-circuit defect can be repaired by means of the so-called passivation treatment. However, the photovoltaic element thus repaired is potentially in a state that it is likely to be again shunted. This situation is facilitated particularly when moisture is present therein.\nIn the case of a solar cell module in which such repaired photovoltaic element is used having a fluororesin film as the surface protective member, the fluororesin film as the surface protective member has difficultly functioning as a barrier to sufficiently prevent moisture from invading into the solar cell module. In this case, when the photovoltaic element is sealed by a filler comprising an EVA having a high hygroscopic property, the solar cell module is not always satisfactory in terms of assuring the stability of the solar cell module upon continuous use under an environment with high temperature and high humidity over a long period of time.\nIn the case where a glass fiber is provided in the filler in order to protect the surface of the photovoltaic element, the above moisture invasion phenomenon is facilitated by way of the capillary action of the glass fiber.\nFor the EVA as filler by which the photovoltaic element is sealed, there is a tendency that the acetic acid radical is hydrolyzed in the present of moisture to liberate acetic acid. The liberated acetic acid is liable to corrode the surface of the transparent and electrically conductive layer, resulting in facilitating to shunt the photovoltaic element."}
{"text": "This invention relates to a vehicle including a loading implement, and more particularly, to a vehicle incorporating a bucket which may be moved in a slow dump and a fast dump mode.\nU.S. Patent Application No. 775,062, entitled \"HIGH PRESSURE IMPLEMENT CIRCUIT FOR LOADER WITH SLOW AND FAST DUMP POSITION\", filed Mar. 7, 1977 (assigned to the assignee of this invention), discloses a system for use with a vehicle including an implement movable to a rack-back and a dump position. As disclosed therein, the operator of the vehicle may select a relatively high speed of movement from the rack-back to the dump position, or a relatively low speed of movement thereof. In such system, it will be seen that once a certain dump mode is chosen, that dump mode will remain in effect until the operator of the vehicle selects a different dump mode. In order to add to the overall flexibility of operation of the apparatus, it would be useful to provide means wherein, upon selection of the fast-dump mode with the implement in the rack-back position, the implement automatically moves to the slow-dump mode after a certain amount of movement thereof, to provide highly effective control of the load carried by the implement."}
{"text": "The present invention relates generally to a system for improving the equalization of overall utilization of wide abrasive belts utilized in surface abrading machines, and more particularly to a use prompter/monitor system for use in combination with such machines wherein an array of individual monitors and/or sensors is employed. Each sensor monitors one discrete segment of the transverse width of the belt, with the activity of each segment being compared to all others in the array. The system of the present invention serves to monitor the extent of usage across the entire width of the belt, and at the same time, serves to prompt the operator when certain segments are being overly utilized or alternatively are being under-utilized. Since the lifetime of a wide abrasive belt relates directly to the most heavily utilized segment, it serves an economic advantage as well as an environmental benefit to equalize belt usage as much as is practicable.\nIn the operation of many machines, and particularly wide-belt sanders, the recommended and preferred practice is to endeavor to have the workpieces or product being surfaced be substantially evenly fed across the working width of the machine. Because of variations in the application of these machines to a variety of different workpieces and due to individual operator preferences, it is generally acceptable if the product being surfaced is more or less evenly fed across the width of the machine.\nAs indicated above, the practice of feeding workpieces evenly is essential in order to distribute wear of machine components in addition to abrasive belts, such as conveyor belts, conveyor belt support beds, hold-down shoes, pinch rolls, contact drum surfaces and the like. It is also desirable to achieve uniform wear of other surface treating or surfacing components such as knife cutter heads, brushes, buffs, or other finishing or abrading-type tooling used on these machines. The term xe2x80x9cabrasive beltsxe2x80x9d as employed herein is utilized in a comprehensive sense, and is intended to include those other moving surface treating components including knife cutter heads, brushes, buffs, or similar abrading-type tooling used on these machines.\nIt has been a long term and widely recognized problem throughout the industry that the seemingly simple task of evenly feeding workpieces or products across the width of a machine is difficult to achieve, and in certain instances, even more difficult to manage. Because of the nature of the individual operations or tasks, machine operators generally prefer to feed products in-line and management personnel and/or supervisors have no practical way to monitor this function so as to ascertain that even or uniform feeding of workpieces is being achieved. The width utilization system of the present invention will not only prompt an operator to evenly feed workpieces across the width of a machine, but will additionally provide supervisory personnel with the means to monitor individual operator performance as it relates to workpiece feed practices.\nBriefly, and in accordance with the present invention, a wide belt utilization prompter/monitor system is provided for use in combination with wide-belt abrasive surface treating apparatus or systems. Wide-belt abrasive surface treating apparatus typically are provided with a frame which supports a conveyor with drive means for moving an endless belt in a drive direction along a defined path, with the belt accordingly having a top flight for transporting individual workpieces along a drive axis between an infeed and a discharge end. One or more workpiece treating stations are operationally positioned between the ends, with the workpiece treating station defining a wide abrading station. Typically, a wide abrasive belt or similar surface treating mechanism is disposed in one or more workpiece treating station with other devices such as, for example, rotary knife cutter heads, brushes, buffs, or other surface finishers being provided as required. Each workpiece treating station is superimposed adjacent the top flight and positioned to contact the surface of workpieces being transported along the conveyor top flight.\nThe abrasive use monitor system of the present invention comprises an array consisting of a plurality of axially spaced apart signal generating workpiece sensors, with these sensors being disposed in axially spaced apart relationship across the width of the belt, preferably transversely to the drive axis. Each sensor creates a signal in response to the presence of a workpiece moving along a certain discrete or predetermined transverse segment of the conveyor belt. A readout means such as a data processor system is provided for receiving signals from each sensor in the array, with the readout means including a readout indicator or information display. The data processor operatively responds to the signals generated by the sensors and the resultant information is displayed on an information display panel. The data processing means are coupled between the sensors and the readout means or information panel, and evaluates and determines the accumulative duration of time during which the presence of a workpiece is detected by each of the individual sensors comprising the array. By accumulating this information, the data processing means determines which individual sensor in the array has detected the shortest or lowest cumulative amount of time. Thus, the data processing means is adapted to compare the duration of cumulative times for which workpieces were detected by each of the individual sensors in the array relative to that certain sensor detecting the shortest duration of cumulative time. As an alternative, the data processing means may be adapted to detect which sensor in the array has the greatest duration of cumulative time of workpiece exposure. These extremes of cumulative exposure may be referred to as the xe2x80x9cthe range limitxe2x80x9d, and hence the term xe2x80x9crange limitxe2x80x9d is used in a comprehensive sense, and is intended to define one or the other of the extremes, either the greatest duration or the least duration. Additionally, the readout means are designed to indicate the differences in cumulative times for each sensor in the array relative to that certain sensor representing the extreme range limit of time, such as the shortest duration of cumulative time.\nAs an added feature of the present invention, the data processing means may include an input representing the conveyor rate of speed for determining the aggregate length of the individual workpieces which have been transported along the top flight and detected by the sensors over any desired period or interval of time. This feature advantageously provides information useful in determining when an individual abrasive belt or surface contacting member such as a knife cutter head is reasonably in need of replacement.\nTherefore, it is a primary object of the present invention to provide an improved system for monitoring the operation of a workpiece surfacing component such as a wide abrasive belt so as to enhance uniform distribution of wear across the transverse width of the surfacing component.\nIt is yet a further object of the present invention to enhance the uniform utilization of a wide surface treating system such as an abrasive belt system wherein an array of individual axially spaced apart sensors are utilized to sense workpiece presence and thereby indicate the extent to which certain width segments of the apparatus have been exposed to contact with workpieces and to assist in uniform distribution of wear.\nIt is a yet a further object of the present invention to provide an improved monitor system for use with wide-belt abrasive surface treating apparatus wherein an array of individual axially spaced apart workpiece sensors are provided to gather data relating to the history of distribution of workpieces across the width of a conveyor-fed wide abrasive surface treating apparatus.\nOther and further objects of the present invention will become apparent to those skilled in the art upon a study of the following specification, appended claims, and accompanying drawings."}
{"text": "This invention relates generally to the field of syringes and, more particularly, to syringes used with viscoelastic agents.\nDuring surgery, particularly ophthalmic surgery, various viscoelastic agents may be introduced into the surgical site. These agents generally are expressed into the surgical site out of a syringe and through a relatively thin cannula. The pressure at the syringe/cannula fitting can be very high due to the high viscosity of the viscoelastic agents. As a result, attempts to express the agent out of the cannula can cause the cannula to become disconnected from the syringe.\nPrior art syringes intended for use with viscoelastic agents, such as described in U.S. Pat. No. 4,540,405 (Miller, et al.) have been directed primarily at providing a better grip on the relatively small syringe so that sufficient force can be applied to express the viscoelastic agent, and the cannula may still become disconnected from the syringe.\nAccordingly, a need continues to exist for a holder that tightly secures a cannula on the end of a syringe."}
{"text": "This invention relates to an air bag restraint system, and more particularly to an air bag restraint system which controls the expansion characteristics of an inflatable air bag cushion discharged therefrom such that the inflated profile of the discharged air bag may be varied based upon the size and/or position of the vehicle occupant to be protected and/or upon the severity of a crash event.\nIt is well known in the prior art to provide an air bag assembly including an inflatable air bag for protecting the occupants of a transportation vehicle. In an automotive vehicle such air bag assemblies are typically located within the hub of the steering wheel and in a recess in the vehicle instrument panel for protection of the vehicle occupants seated in opposing relation to such assemblies. Additional air bag assemblies may be located within the seats and/or door panels for protection of the occupants during a side-impact event. It is also known to utilize inflatable curtain-like structures for deployment from the structural pillars or roof line of the motor vehicle so as to promote restraint and protection of the vehicle occupant during a roll-over event.\nAir bag assemblies typically include an inflatable cushion in fluid communication with a gas emitting inflator. Upon sensing certain predetermined vehicle conditions, such as a certain amount of vehicle deceleration, the inflator discharges a fixed amount of inflator gas thereby forcing the air bag into a deployed position. The inflator gas occupies the available volume within the air bag cushion thereby forcing the air bag cushion to expand outwardly to the extent permitted by its construction. The pressure within the air bag cushion upon deployment is proportional to the quantity of inflator gas expelled into the air bag and inversely proportional to the volume occupied by the inflator gas within the air bag. As the occupant comes into contact with the expanded air bag, the inflator gas is forced out of the air bag thereby dissipating the kinetic energy of the occupant.\nIn some cases, it may be desirable to provide an inflator that has varied levels or stages of inflator gas output in response to the sensing of different vehicle or occupant conditions. Thus, it is generally known in the prior art to provide multi-stage inflators that discharge inflation gas at variable levels depending upon the conditions present during deployment. However, these multi-stage inflators are more complex than typical inflators. Moreover, the use of such multi-stage inflators provides control over only the amount of inflator gas which is discharged and does not provide control over the expanded geometry of the inflated air bag cushion. That is, so long as the air bag has a fixed expanded geometry, the inflator gas will tend to fill the available capacity and the expanded configuration of the air bag will be generally the same even if the quantity of inflator gas is varied although the pressure within the air bag will tend to differ appreciably.\nIn order to provide an additional degree of freedom in the control of air bag performance, it has been suggested to utilize air bag cushions which incorporate seams within the air bag to control the expanded geometry of the inflated air bag wherein the seams separate upon the introduction of sufficient force across the seams thereby freeing the air bag cushion from the restraint imposed by the seams at lower pressures. In order for such break-away seams to provide controlled expansion, the application of such seams must be performed with substantial precision such that seam separation will occur in a highly reproducible and predictable manner. As will be appreciated, due to the large number of variables involved in the introduction and separation of such break-away seams, such requisite precision and reproducibility may be difficult to achieve.\nThe present invention provides advantages and alternatives over the prior art by providing an assembly to vary the deployed profile of an air bag cushion which may be inflated to a controlled variable level using either a single stage or multi-stage inflator and which may be expanded preferentially in the depth direction towards the occupant to be protected without the need to rely exclusively on break-away seams. The present invention thereby provides an air bag cushion which may be deployed to a relatively shallow profile during a crash event of low severity or for the protection of a small stature occupant or an occupant seated in close proximity to the air bag and which may also be deployed to a profile of substantially increased depth during a severe crash event or for protection of a larger occupant or an occupant seated further away from the location of air bag deployment. The present invention further provides a mechanism for simply and effectively controlling both the profile of the deployed air bag cushion and the quantity of inflation gas released into the air bag cushion.\nAdvantageously, according to one aspect this invention provides an assembly to deploy an air bag cushion to a controlled geometry of a depth which is desired in view of the position of the occupant to be protected such that an air bag of substantial depth is available to protect occupants such as larger stature persons who are seated further away from the location of cushion deployment, while a smaller air bag volume of diminished depth is available to protect occupants such as smaller stature persons seated closer to the location of cushion deployment. The deployed profile may also be adjusted based upon the severity of the crash event giving rise to the deployment.\nThe present invention further provides an effective mechanical actuation mechanism to adjust the deployment characteristics of the air bag cushion which actuation mechanism is not dependent upon a particular cushion construction technique and is activated based upon measured parameters of the crash event and the occupants to be protected.\nFurthermore, the actuation mechanism to control the deployed profile of the air bag cushion may be operated either independently or in conjunction with a device to adjust the level of inflator gas which enters the air bag cushion thereby providing additional versatility of operation. Conjunctive operation may be carried out utilizing a single common initiator such as a pyrotechnic squib. Thus, it is a potentially preferred feature of the present invention that the air bag assembly may incorporate as few as two initiators (i.e. one initiator to activate the inflator and one initiator to initiate the actuation mechanism to control the deployment profile of the air bag and the level of inflation gas which enters the air bag).\nIn addition, the actuation mechanism to control the deployed profile and level of inflation gas may be mounted externally on the inflator or the supporting housing thereby avoiding any need to substantially modify the shape or construction of the inflator housing.\nIn addition, the actuation mechanism and variable profile air bag of the present invention may be used in conjunction with traditional single stage inflators without requiring internal modifications to such inflators.\nThese advantages are accomplished in a potentially preferred form of the present invention by providing an inflatable air bag expandable to a first volume and depth which first volume and depth may be restrained to a second diminished volume and depth by tethers which are attached at locations across the surface of the air bag cushion and which are releasably attached to a releasable anchor assembly independent from the air bag cushion. Based upon the severity of the crash event and/or the size and/or the position of the occupant to be protected, an actuation mechanism which releases the tethers from the anchor assembly may be either activated or may remain deactivated to yield a deployed air bag cushion of desired volume and geometric configuration. The volume of inflation gas entering the air bag cushion may be adjusted simultaneously with the profile of the air bag cushion.\nAccording to the potentially preferred form of the invention, the actuation mechanism will remain deactivated in the event that the crash event is of low severity and/or the occupant to be protected is of small stature and/or the occupant to be protected is seated in a position close to the location of air bag deployment, thereby resulting in a constrained final expanded geometry. In the event that the crash event is of sufficient severity and/or the occupant to be protected is of larger stature and/or the occupant to be protected is seated further away from the location of air bag deployment, the actuation mechanism may be activated to at least partially eliminate depth restraint thereby permitting the air bag to expand to a greater volume and depth so as to provide additional cushioning for such an occupant.\nThus, it will be appreciated that variable expanded geometries may be achieved in the inflated air bag with such geometries being controllable and adjustable based upon the severity of the crash event and/or the physical parameters of the occupant to be protected. Such control is preferably achieved through use of one or more mechanical actuation mechanisms which are either activated or deactivated to effect control of the inflatable air bags by either holding or releasing tethering restraint elements from an anchoring location remote from the air bag cushion.\nSince the activation or deactivation of the actuation mechanism is carried out in conjunction with deployment of the air bag cushion, the necessary geometric character may be achieved for a given occupant as may be present within the vehicle at the time of deployment. Thus, the variable air bag and actuation mechanism according to the present invention offers substantial versatility in the protection of large numbers of diverse occupants including very small and very large occupants experiencing different levels of crash severity. It will also be appreciated that this system provides a simple, cost effective and highly reproducible system for controlling the volume of deployed air bag cushions even when using a traditional single stage inflator."}
{"text": "Conventionally, a typical driving system of inverter is a Royer's circuit (for example, see the patent literature 1). FIG. 8 shows an example of the basic circuit of Royer's circuit.\nIn FIG. 8, a transformer 1 has two primary windings L1, L2 and a single secondary winding L3, and a negative electrode of DC power source 2 is connected to a mutually connected end of the primary windings L1, L2, and a positive electrode of DC power source 2 is connected to the other end of the primary winding L1 via a switch SW1 while being connected to the other end of the primary winding L2 via a switch SW2.\nAs shown in FIG. 9(a), the switch SW1 and the switch SW2, both of which are initially turned off, are alternately turned on and off during the time T1, T2, allowing to alternately flow positive and negative input current in the primary windings L1, L2 of transformer 1, whereby an output signal as shown in FIG. 9(b) can be obtained in the secondary winding L3 of transformer 1.\nAccording to the driving system of this Royer's circuit, even if the input current may instantaneously become zero during operation, the operation of circuit is not completely suspended.\nFurther, another typical driving system of inverter is PWM (Pulse Width Modulation) (see FIGS. 10(a), 10(b)). This is a system in which the pulse widths T1, T2 of input current in the cycle T of on-and-off operations of switches SW1, SW2 is adjusted to allow the intermittent driving of circuit, whereby the secondary output voltage is controlled.\nThat is, the pulse width of the secondary output voltage is changed by adjusting T1 and T2 which are on-time of the switches SW1 and SW2 in FIG. 8, thereby adjusting an effective output voltage. Thus, if the input current of transformer 1 is within the range of increasing in proportion to the time T1, T2 and the magnetic field of transformer 1 is within the range of saturation, the longer the time T1, T2 (that is, the higher the duty), the greater the effective value of output voltage becomes.\nHereby, PWM is a system in which input current is intermittently driven for the purpose of “control of output voltage”."}
{"text": "This invention relates to portable vacuum cleaners and, more particularly, to a vacuum cleaner capable of operating both with air and liquid cleaners.\nVacuum cleaners are utilized in numerous situations ranging from relatively light duty, such as the removal of crumbs and dust from a flat surface, as well as for relatively heavy duty operation as in the withdrawal of foreign matter embedded in carpets and upholstery. It is apparent from the wide range of cleaning tasks that some cleaning is best accomplished by the use of air alone, while other cleaning is best accomplished with the use of water or other cleaning liquid which is to be drawn by suction into the vacuum cleaner.\nIt is recognized that the use of a liquid cleaning agent necessitates a more complex structure in the vacuum cleaner. Thus, it is necessary to protect a fan motor from contamination by the liquid. Provision must also be made for extraction of the foreign matter and the liquid cleaner from the vacuum cleaner upon completion of the cleaning process. In addition, the foregoing must be accomplished while allowing for the intake and exhaust of the air stream which is driven by suction of the cleaner fan.\nThe foregoing constraints become more difficult to attain in the case of a portable vacuum cleaner, since, as is readily appreciated, a hand held cleaner may be placed in a variety of positions and orientations so that, unlike a stationery cleaner, reliance cannot be made solely on the use of gravity for direction of the liquid cleaning agent away from the motor. Also, it is realized that the use of the traditional vacuum-cleaner bag fabricated of cloth or paper would be contraindicated since any liquid entrapped therein would tend to leak out upon removal of the bag.\nThus, a problem exists in that the desirable feature of portability in a vacuum cleaner is difficult to attain in a situation wherein the vacuum cleaner is to be used for both wet and dry cleaning applications."}
{"text": "Wireless communications systems, such as the third-generation (3G) of mobile telephone standards and technology are well known. Such 3G standards and technology have been developed by the Third Generation Partnership Project (3GPP). The 3rd generation of wireless communications has generally been developed to support macro-cell mobile phone communications. Communication systems and networks have developed towards a broadband and mobile system.\nThe 3rd Generation Partnership Project has developed the so-called Long Term Evolution (LTE) system, namely, an Evolved Universal Mobile Telecommunication System Territorial Radio Access Network, (E-UTRAN), for a mobile access network where one or more macro-cells are supported by a base station known as an eNodeB or eNB (evolved NodeB). More recently, LTE is evolving further towards the so-called 5G or NR (new radio) systems where one or more cells are supported by a base station known as a gNB (generalized NodeB).\nIn 3GPP there is a continual search for agreements on the implementation of new and relevant features for NR systems. The agreements are relevant to many options. One option of interest relates to the use of Polar coding for the eMBB (enhanced Mobile Broadband) UL/DL control channels. According to this option, CRC (cyclic redundancy check) codes will be utilized either within the coded block or external in another coded block in the case of concatenated polar codes. These CRC codes may be used to provide error detection for the blocks of data being transmitted and, in some implementations may also enable error correction.\nPolar codes in 5G have been proven to achieve capacity by code construction for binary erasure channel (BEC). Encoding complexity is of O(Nlog2N) for a code block size N. Decoding complexity is of O(Nlog2N) with a successive cancellation (SC) decoder, and higher with more advanced decoder with better performance such as a List decoder O(LNlog2n). The implementation complexity of list decoder increases with increasing list size, especially with large block size. Furthermore, Polar codes are not parallelizable as Turbo and LDPC, and latency requirements for large block-lengths are questionable. However, Polar codes have been shown to outperform Turbo code and LDPC for short-block lengths with no error floor for List decoder, as NR control channel requires ultrareliability transmission.\nIn a communications system, encoding is generally used to improve the reliability of data transmission. The polar code is a linear block code which can provably achieve the Shannon capacity for a binary erasure channel with low complexity of encoding-decoding. Conceptually, polar codes split a physical channel up into a number of virtual channels. The reliability of each of the virtual channels (i.e. the likelihood of a bit transmitted on that channel will be successfully received) defined by a polar code differs from each other. However, generally, as the number of virtual channels defined by a polar code increases (e.g. as N goes to infinity), the reliability of each of the channels polarizes to become either very reliable (e.g. a perfect channel with capacity 1) or very unreliable (e.g. a useless channel with capacity 0). The virtual channels with the highest reliability are chosen for the transmission of information bits, whilst the other virtual channels are “frozen” and are not used for the transmission of information bits. The positions of data and frozen bits change with the polar code construction parameter  for the same code block size. Both the encoder and the decoder must utilize same locations of the frozen bits for a successful decoding meaning they should know either parameter  for code construction or a pre-defined set of frozen locations which is the outcome of the code construction.\nAs illustrated in FIG. 1, a polar code 120 of size N (whereby N=8 in FIG. 1) is generated by duplicating and combining n times the constructed transform starting from duplicating and combining the basic transform 110, to give N virtual channels, where N=2n. It should be noted that the symbol ⊕, as used in FIG. 1, represents the exclusive OR (XOR) operation. Of these channels, the four channels having the highest reliability are indicated as being “data” channels, whilst the four channels having the lowest reliability are indicated as being “frozen” channels, for a code rate R=½.\nVarious agreements have been made in this technical area, for example:—\nAgreement:\n                J CRC bits are provided (which may be used for error detection and may also be used to assist decoding and potentially for early termination)                    J may be different in DL and UL            J may depend on the payload size in the UL (0 not precluded)                        In addition, J′ assistance bits are provided in reliable locations (which may be used to assist decoding and potentially for early termination)        J+J′<=the number of bits required to satisfy the FAR target (nFAR)+6                    Working assumption:                            For DL, nFAR=16 (at least for eMBB-related DCI)                For UL, nFAR=8 or 16 (at least for eMBB-related UCI;                                    note that this applies for UL cases with CRC)            J′>0            Working assumption: J″<=2 additional assistance bits are provided in unreliable locations (which may be used to assist decoding and potentially for early termination)                            Can be revisited in RAN1#89 if significant benefit is shown from a larger value of J″ without undue complexity—companies are encouraged to additionally evaluate J″=8                                    The J′ (and J″ if any) bits may be CRC and/or PC and/or hash bits (downscope if possible)            Placement of the J, J′ (and J″ if any) assistance bits is FFS after the study of early termination techniques                            Appended?                Distributed?                evenly?                unevenly?Agreements (for very small block lengths):                                                K=1 (if channel coding is applied):                    Repetition code                        K=2 (if channel coding is applied):                    Simplex code                        3<=K<=11:                    LTE RM code                            Note that if NR requires a codeword size N that is not supported by the LTE RM code, then the LTE RM code will be extended by repetition as in LTE                                                12<=K:                    Polar code (single design for all control information sizes, except for possible omission of CRC bits for payloads <=−22 bits)Agreement for DCI:                        Maximum mother code size of Polar code, N=2n, is:                    Nmax,DCI=512 for downlink control informationWorking Assumption for UCI:                        Nmax,UCI=1024                    Optimise code design for K up to 200                            Also aim for code design that supports values of K up to 500 with good performance, typically using higher code rates                                                Without prejudice to the final design, companies are encouraged to investigate advanced code rate matching schemes until RAN1#88bis        Working assumption can be revisited at RAN1#88bis if it does not prove to be possible to generate a good code design with Nmax, UCI=1024Agreements:        Performance metrics (may be based on analytic derivation)                    BLER            FAR (with AWGN as input to the decoder)                        Polar codes for control channels support one of the following alternatives:                    Alt. 1: CRC+“basic polar” (i.e. as per above agreed description) codes                            1a: Longer CRC                                                e.g. (J+J′) bits CRC+basic polar                    1b: J bit CRC                            The J bits can be distributed                                    The CRC can be used for both error detection and error correction                        Alt. 2: J bits CRC+concatenated polar codes                    e.g. J bits CRC+J′ bits CRC+basic polar;                            J bits CRC+J′ bits distributed CRC+basic polar;                J bits CRC+PC bits+basic polar; (i.e. PC-Polar)                J bits CRC+Hash sequence+basic polar;                . . .                                    J bits CRC is only used for error detection                        \nPolar coding is adopted for eMBB UL/DL control channels (LPDC was adopted as the coding scheme for eMBB UL/DL data channels). Coding scheme(s) for URLLC and mMTC is not yet defined.\nPolar code will utilize CRC, either within the coded block or externally in another coded block in case of concatenated polar codes."}
{"text": "Packet-based data networks continue to grow in importance, and it is often desirable to monitor network traffic associated with these packet-based networks on an ongoing basis. To meet these monitoring needs, copies of network packets can be forwarded to diagnostic network monitoring tools. Packets are often forwarded using network hubs, test access ports (TAPs), and/or switched port analyzer (SPAN) ports available on network switches. For example, certain network switches produced by Cisco Systems include SPAN ports to which traffic on the switches are mirrored. It is also noted that other packet monitoring or access methods may also be used to acquire copies of network packets being communicated within a network infrastructure.\nTo help alleviate the problem of limited access to network packets for monitoring, tool aggregation devices have been developed that allow shared access to the monitored network packets. These tool aggregation devices allow users to obtain packets from one or more network monitoring points (e.g., network hub, TAP, SPAN port, etc.) and to forward them to different monitoring tools. U.S. Pat. Nos. 8,018,943 and 8,098,677 describe example embodiments for network tool optimizers that provide solutions for packet filtering and provide, in part, configuration of user-define filters, automatic creation of filter engine forwarding rules, automatic handling of filter overlaps, graphical user interfaces (GUIs) for filter creation, and other features. U.S. Pat. Nos. 8,018,943 and 8,098,677 is each hereby incorporated by reference in its entirety.\nWhen a network to be monitored includes a cloud or virtual environment, however, difficulties arise in utilizing prior tool aggregation devices particularly where multiple unrelated users are allocated different cloud resources within the cloud or virtual environment. One such example environment is where resources within a server cloud are offered by a controlling entity (e.g., Amazon Web Services) to different user entities that lease, rent, or otherwise pay for server cloud resources from the controlling entity. If these user entities desire to monitor network activity, the preferred place to conduct such monitoring is often within the network infrastructure for the server cloud. However, as the server cloud is not controlled by or directly accessible to the user entities, this cloud network monitoring is impractical to implement using current tool aggregation devices without the controlling entity opening its network infrastructure to the user entities. Further, the controlling entity of the server cloud typically does not want its network infrastructure or traffic to be visible to the user entities\nEven assuming the controlling entity (e.g., Amazon Web Services) is willing to connect a monitoring tool to its server cloud so that copies of packet traffic are forwarded to a monitoring tool (e.g., through a hub, TAP, SPAN port, etc.) desired to be used by a user, these packets will likely include packets for other user entities in addition to the user desiring the network monitoring. This overlap in packet traffic is likely to occur because cloud resources are often virtualized such that processors or cores associated with any particular physical server processing platform (e.g., server blade) may be used by two or more user entities through server instances created within the server processing system platform. Traffic copied from this physical server processing system platform, therefore, will likely include packets that are associated with multiple independent user entities using the different server instances. As such, significant security and confidentiality issues arise if an attempt is made to add network monitoring tools and associated monitoring services to server cloud resources being used and shared by user entities."}
{"text": "1. Technical Field\nApparatuses and methods consistent with exemplary embodiments relate to an X-ray imaging apparatus and a method of controlling the same, and more specifically, to an X-ray imaging apparatus that reconstructs a current frame using previous frame information and a method of controlling the same.\n2. Description of the Related Art\nAn X-ray imaging apparatus is an apparatus capable of obtaining an internal image of an object by radiating X-rays onto the object and collecting the X-rays transmitted through the object. Because the permeability of objects, materials, and/or elements within the object differ with regards to allowing X-rays to pass through depending on properties of substances composing the object, it is possible to image an internal structure of the object by detecting an intensity or a strength of X-rays transmitted through the object.\nIn order to provide a level of safety from the X-ray imaging apparatus, reducing a dose of X-rays incident on the object is recognized as an important issue and thus a great deal of research and development for reducing an X-ray dose is underway. One technological implementation for reducing an X-ray dose includes a method of radiating X-rays onto only a region of interest which has been generally adopted."}
{"text": "1. Field of the Invention\nThe present invention relates to a rolling element accommodating belt, a linear guide apparatus and a metallic mold for manufacturing the rolling element accommodating belt.\n2. Description of Related Art\nIn the linear guide apparatus, a slider is relatively moved with respect to a guide rail through a plurality of rolling elements which are rolling and circulating in an infinite circulating passage. However, in the linear guide apparatus, while the slider is relatively moving with respect to the guide rail, the rolling elements are rolling and moved in the same direction. Therefore, the rolling elements, which are adjacent to each other, are rubbed to each other. Accordingly, the rolling elements can not be smoothly rolled. Therefore, an intensity of noise is increased and abrasion of the rolling elements is facilitated. In order to solve the above problems by suppressing the generation of noise and making the linear guide apparatus operated smoothly, a rolling element accommodating belt is conventionally proposed in which the rolling elements are aligned in the alignment direction in the infinite circulating passage so as to smoothly operate the linear guide apparatus. Concerning this rolling element accommodating belt, for example, refer to Japanese Patent Unexamined Publications JP-A-05-52217, JP-A-2001-165169, JP-A-10-9264, JP-A-09-14264, JP-A-11-247856 and JP-A-2005-69444.\nFor example, according to the JP-A-05-52217, JP-A-2001-165169 and JP-A-10-9264, a rolling element accommodating belt is disclosed which includes spacer portions interposed between the rolling elements adjacent to each other and also includes connecting arm portions for connecting the spacer portions. According to the rolling element accommodating belt composed as described above, the rolling elements are aligned in a line in the alignment direction in the infinite circulating passage so as to suppress the generation of noise so that the rolling elements can be smoothly circulated in the infinite circulating passage. In this connection, according to the technique described in the JP-A-2001-165169, the rolling element accommodating belt is composed in such a manner that the rolling elements accommodated on the rolling element accommodating belt can be freely detached in a perpendicular direction to the front and the back surface of the rolling element accommodating belt.\nThis type rolling element accommodating belt is obtained when melted resin material is poured from a gate into a product configuration portion in a metallic mold. However, in the product configuration portion in the metallic mold, residual air and gas, which is generated from the melted resin material, exists. These gases tend to stay in a portion where a gas flowing passage comes to a dead end. Especially, these gases tend to stay in an end portion of the above spacer portion. Therefore, it is difficult for the melted resin to smoothly flow into this portion where the gases are staying. Accordingly, an underfill portion of resin (which is a portion where predetermined amount of resin is not filled) is generated on the rolling element accommodating belt formed in the mold. When the underfill portion of resin is generated in this way, there is a possibility that damage starts from this underfill portion of resin generated on the rolling element accommodating belt.\nIn this case, for example, according to the technique described in the JP-A-11-247856 or JP-A-2005-69444, a joint portion of the upper mold and the lower mold of the metallic mold can be made to function as a gas vent. However, at the end portion of the formed spacer portion, no means is provided for discharging gas. Therefore, gas tends to stay at the end portion of the formed spacer portion. Therefore, the spacer portion needs to be more investigated to discharge gas so that the generation of an under fill portion of resin can be prevented.\nFurther, concerning the method of manufacturing this type rolling element accommodating belt, for example, the JP-A-05-52217 discloses a technique in which the rolling elements are arranged in a metallic mold so that the rolling elements can be used as a core and then injection molding is conducted to manufacture the rolling element accommodating belt.\nThe JP-A-09-14264 discloses a technique in which injection molding is conducted while the rolling elements are being used as a core in the same manner as that described above. However, in order to solve a problem that a portion, in which the rolling elements are accommodated, and the rolling elements are tightly contacted with each other by the shrinkage caused at the time of forming, oil or water is absorbed after the completion of molding.\nAnother manufacturing method is disclosed in the JP-A-11-247856 as follows. For example, in the JP-A-11-247856, rolling element shaped-molds, the sizes of which are larger than the diameters of the rolling elements to be used, are arranged at predetermined intervals and injection molding is conducted to manufacture a rolling element accommodating belt.\nThe JP-A-2005-69444 discloses a technique in which a formed rolling element accommodating belt is detached from a metallic mold by obliquely moving the metallic mold used in the injection molding. In this case, in the examples disclosed in the JP-A-11-247856 and JP-A-2005-69444, a position of a dividing line of an upper mold and a lower mold is located at a peripheral portion of a spacer portion. Concerning this matter, refer to FIG. 2 of the JP-A-11-247856 and FIG. 9 of JP-A-2005-69444.\nHowever, in the above method of manufacturing the rolling element accommodating belt, for example, in the technique described in the JP-A-05-52217 in which injection molding is conducted while the rolling elements are being used as a core, due to the shrinkage of material caused at the time of forming, sizes of portions, in which the rolling elements are accommodated, are reduced smaller than the predetermined sizes. Therefore, the rolling elements can not be smoothly rolled.\nFor example, in the technique described in the JP-A-09-14264, in order to avoid a case where the rolling elements can not be smoothly rolled because of the shrinkage caused at the time of forming, the processing of absorbing oil or the processing of absorbing water is conducted after the completion of forming. In this case, it is necessary to strictly control the conditions of processing of absorbing oil or absorbing water, which raises the manufacturing cost.\nFor example, in the technique described in the JP-A-11-247856, in the spacer portion which has been formed, a concave face corresponding to the circumferential face of the rolling element mold is formed as an under-cut portion. Therefore, when the rolling element belt is detached from the metallic mold, a strong force is given to the rolling element accommodating belt. Accordingly, there is a possibility that the rolling element accommodating belt is damaged or deformed.\nFurther, in the technique described in the JP-A-11-247856 and the JP-A-2005-69444, a position of the dividing line of the upper and the lower mold of the metallic mold is located in the peripheral portion of the spacer portion. Accordingly, burr is generated in the peripheral portion of the spacer portion. Due to the generation of burr, there is a possibility that the linear guide apparatus can not be smoothly operated. That is, at the time of forming, resin flows into a joint portion of the upper mold and the lower mold, which causes the generation of burr. In the case where this burr is generated in the peripheral portion of the spacer portion, when the rolling element accommodating belt circulates in the infinite circulating passage of the linear guide apparatus, the rolling element accommodating belt rubs an inner circumferential wall of the infinite circulating passage. Alternatively, the rolling element accommodating belt is hooked on the inner circumferential wall of the infinite circulating passage. Therefore, there is a possibility that a smooth operation of the linear guide apparatus is obstructed.\nIn the case of the rolling element accommodating belt disclosed in JP-A-2001-165169, the rolling elements accommodated there can be freely detached in a direction perpendicular to the surface side and the back side of the rolling element accommodating belt. For example, as exemplarily shown in FIG. 17, a position of the dividing line of the metallic mold of the spacer portion is formed while crossing a face including a portion coming into contact with the rolling elements. When the dividing line of the metallic mold is set at this position, burr is not generated in the periphery of the spacer portion, however, burr is generated in a portion coming into contact with the rolling elements. Therefore, by the burr protruding onto the rolling element side, the face, which is to be contacted with the rolling elements, and the rolling elements can not be stably contacted with each other. Accordingly, smooth operation of the linear guide apparatus is obstructed. In order to prevent the above generation of burr, it is necessary to enhance the accuracy of the joint face of the metallic mold. Alternatively, it is necessary to severely control the forming condition. However, when the above countermeasures are taken, the manufacturing cost is raised."}
{"text": "Many human enzymes serve as targets for the action of pharmaceutically active compounds. Several classes of human enzymes that serve as such targets include helicase, steroid esterase and sulfatase, convertase, synthase, dehydrogenase, monoxygenase, transferase, kinase, glutanase, decarboxylase, isomerase and reductase. It is therefore important in developing new pharmaceutical compounds to identify target enzyme proteins that can be put into high-throughput screening formats. The present invention advances the state of the art by providing novel human drug target enzymes related to the aminoacylase subfamily.\nThe present invention has a substantial similarity to aminoacylase-1. Aminoacylase-1 (ACY1, EC 3.5.1.14), a new type of metalloprotein, is a cytosolic enzyme with a wide range of tissue expression and has been postulated to function in the catabolism and salvage of acylated amino acids. ACY-1 is more highly expressed in kidney than in liver. ACY1 has been assigned to chromosome 3p21, a region reduced to homozygosity in small-cell lung cancer and renal cell carcinoma, and shows a reduced or absent expression in small-cell lung cancer cell lines and tumors. For a review related to aminoacylase-1, see Miller et al., Genomics 1990 Sep.; 8(1):149-54, Mitta et al., J Biochem (Tokyo) 1992 Dec.; 112(6):737-42.\nEnzyme proteins, particularly members of the aminoacylase subfamily, are a major target for drug action and development. Accordingly, it is valuable to the field of pharmaceutical development to identify and characterize previously unknown members of this subfamily of enzyme proteins. The present invention advances the state of the art by providing previously unidentified human enzyme proteins, and the polynucleotides encoding them, that have homology to members of the aminoacylase subfamily. These novel compositions are useful in the diagnosis, prevention and treatment of biological processes associated with human diseases."}
{"text": "1. Field of the Invention\nThe present invention relates to systems and methods for producing or delivering heat at or near the down hole end of production tubing of a producing oil or gas well for improving production therefrom.\n2. Background Information\nFree-flowing oil is increasingly difficult to find, even in oil wells that once had very good flow. In some cases, good flowing wells simply xe2x80x9cclog upxe2x80x9d with paraffin. In other cases, the oil itself in a given formation is of a viscosity that it simply will not flow (or will flow very slowly) under naturally ambient temperatures.\nBecause the viscosity of oil and paraffin have an inverse relationship to their temperatures, the solution to non-flowing or slow flowing oil wells would seem fairly straight forwardxe2x80x94somehow heat the oil and/or paraffin. However, effectively achieving this objective has proven elusive for many years.\nIn the context of gas wells, another phenomenaxe2x80x94the buildup of iron oxides and other residues that can obstruct the free flow of gas through the perforations, through the tubing, or bothxe2x80x94creates a need for effective down hole heating.\nDown hole heating systems or components for oil and gas wells are known (hereafter, for the sake of brevity, most wells will simply be referred to as xe2x80x9coil wellsxe2x80x9d with the understanding that certain applications will apply equally well to gas wells). In addition, certain treatments (including xe2x80x9chot oil treatmentsxe2x80x9d) for unclogging no-flow or slow-flow oil wells have long been in use. For a variety of reasons, the existing technologies are very much lacking in efficacy and/or long-term reliability.\nThe present invention addresses two primary shortcomings that the inventor has found in conventional approaches to heating oil and paraffin down hole: (1) the heat is not properly focused where it needs to be; and (2) existing down hole heaters fail for lack of design elements which would protect electrical components from chemical or physical attack while in position.\nThe present inventor has discovered that existing down hole heaters inevitably fail because their designers do not take into consideration the intense pressures to which the units will be exposed when installed. Such pressure will force liquids (including highly conductive salt water) past the casings of conventional heating units and cause electrical shorts and corrosion. Designers with whom the present inventor has discussed heater failures have uniformly failed to recognize the root cause of the problemxe2x80x94lack of adequate protection for the heating elements and their electrical connections. The down hole heating unit of the present invention addresses this shortcoming of conventional heating units.\nResearch into the present design also reveals that designers of existing heaters and installations have overlooked crucial features of any effective down hole heater system: (1) it must focus heat in such a way that the production zone of the formation itself is heated; and (2) heat (and with it, effectiveness) must not be lost for failure to insulate heating elements from up hole components which will xe2x80x9cdrawxe2x80x9d heat away from the crucial zones by conduction.\nHowever subtle the distinctions between the present design and those of the prior art might at first appear, actual field applications of the present down hole heating system have yielded oil well flow rate increases which are multiples of those realized through use of presently available down hole heating systems. The monetary motivations for solving slow-flow or no-flow oil well conditions are such that, if modifying existing heating units to achieve the present design were obvious, producers would not have spent millions of dollars on ineffective down hole treatments and heating systems (which they have done), nor lost millions of dollars in production for lack of the solutions to long-felt problems that the present invention provides (which they have also done).\nIt is an object of the present invention to provide an improved down hole heating system for use in conditioning oil and gas wells for increased flow, when such flow is impeded because of viscosity and/or paraffin blockage conditions.\nIt is another object of the present invention to provide an improved design for down hole heating systems which has the effect of more effectively focusing heat where it is most efficacious in improving oil or gas flow in circumstances when such flow is impeded because of oil viscosity and/or paraffin blockage conditions.\nIt is another object of the present invention to provide an improved design for down hole heating systems for oil and gas wells which design renders the heating unit useful for extended periods of time without interruption for costly repairs because of damage or electrical shorting caused by unit invasion by down hole fluids.\nIt is another object of the present invention to provide an improved method for down hole heating of oil and gas wells for increasing flow, when such flow is impeded because of viscosity and/or paraffin blockage conditions.\nIn satisfaction of these and related objects, the present invention provides a down hole heating system for use with oil and gas wells which exhibit less than optimally achievable flow rates because of high oil viscosity and/or blockage by paraffin (or similar meltable petroleum byproducts). The system of the present invention, and the method of use thereof, provides two primary benefits: (1) the involved heating unit is designed to overcome an unrecognized problem which leads to frequent failure of prior art heating unitsxe2x80x94unit invasion by down hole heating units with resulting physical damage and/or electrical shortages; and (2) the system is designed to focus and contain heat in the production zone to promote flow to, and not just within, the production tubing."}
{"text": "1. Field of the Invention\nThe present invention relates to the technical field of reinforcing materials adapted to making up composite parts. More precisely, the invention relates to a novel intermediate material for producing composite parts by subsequent injection or infusion of thermosetting resin, in the form of a veiled tape with an improved resistance to delamination.\n2. Description of Related Art\nComposite parts or articles, i.e. comprising both one or more pieces of reinforcement or fibrous sheets and also a matrix made principally of the thermosetting (resin) type and that may include thermoplastics, may, for example, be produced by using a “direct” or “LCM” (liquid composite molding) method. A direct method is defined by the fact that one or more pieces of fibrous reinforcement are used in the “dry” state (i.e. without the final matrix), the resin or matrix being employed separately, for example by injection into the mold containing the fibrous reinforcement (“RTM” method, resin transfer molding), by infusion through the thickness of the fibrous reinforcement (the “LRI” or liquid resin infusion method, or the “RFI” or resin film infusion method) or by manual coating/impregnation, using a roller or brush, onto each of the individual layers of fibrous reinforcement, applied in succession to the form.\nOther methods known as indirect methods use pre-impregnated materials that already comprise a sufficient quantity of resin to make up the desired composite part. Such materials are in particular those described in documents US 2005/048280, WO 92/20521 and EP 0 554 950.\nFor RTM, LRI, or RFI methods, in general a fibrous preform needs to be produced in the shape of the desired finished article, and then that preform is impregnated with a resin. The resin is injected or infused by pressure or temperature differentials, and then once all of the necessary quantity of resin is contained in the preform, the impregnated preform is heated to a higher temperature in order to carry out the cycle of polymerization/curing and thus cause it to harden.\nThe composite parts used in the automotive, aviation, or shipbuilding industries in particular are governed by very strict regulations, in particular in terms of mechanical properties. It is thus particularly important to have access to materials that are both extremely regular and also easy to handle and use.\nIn those sectors, a large number of preforms are produced that are based on reinforcing materials, formed from carbon fibers, in particular of the unidirectional type. In order to satisfy the high standards in terms of quality and productivity demanded in the aviation sector in particular, it is becoming ever more necessary to use automated methods.\nThe prior art proposes unidirectional sheets of reinforcing yarns in which the cohesiveness between the yarns is ensured by thermoplastic or glass/thermoplastic binding yarns, which may be woven or nonwoven, that extend transversely to the reinforcing yarns. Such sheets are, for example, supplied under the references PW-BUD by SIGMATEX UK Limited, Runcorn, Cheshire WA7 1TE, United Kingdom) or with other denominations (the TeXtreme® range from OXEON (Norrby Langata 4S, SE-50435, Boras, Sweden).\nOther documents such as Japanese patent application JP 2009-235175, proposes depositing a thermosetting resin powder on a fabric type or unidirectional fibrous type support surface in order to provide a base material for producing a preform that has excellent deformability, form stability, good permeability to air, and satisfactory impregnability by resin. Japanese patent application JP 2009 235175 also proposes the same type of intermediate material adapted to the RTM method that can be readily molded and shaped to produce a preform and that can avoid the loss of properties suffered by an interply preform.\nIn order to provide yarns with better cohesiveness, the Applicant has proposed associating each face of a unidirectional reinforcing sheet with a nonwoven that could result in a more continuous bond compared with the spot bonds obtained with binding yarns or powder. Such materials are described in particular in patent application WO 2010/046609. Next, in patent application WO 2010/061114, the Applicant has described a method of producing unidirectional sheets with a given width that have high regularity, adapted to direct methods of producing composite parts from one or more yarns, while limiting material losses. Such a method can in particular be used to obtain controlled width veiled yarns that can then be used to manufacture preforms with a plurality of layers directly using an automated lay-up device. Another route consists in using them in order to make up woven or braided reinforcements that can then be stacked to produce composite preforms or parts using a direct method.\nWhen using the veiled tapes previously proposed by the Applicant, it has been observed that during automated lay-up of a veiled tape, it is bonded to the preceding ply by a combination of a pressure and heating action followed by cooling, where cooling is possibly accomplished without adding a specific coolant, by using a “natural” route. The tape is thus bonded to the preceding ply via its lower face, and that mechanical bond, which is subjected to shear all the time the tape is being laid-up, is of an intensity proportional to the (lay-up tension)/(bonded length) ratio. The lay-up tension is generally assumed to be constant, and so the shear stress is higher during the first centimeters of lay-up and decreases as the length of the laid-up tape increases. The shear force is distributed over the whole thickness of the tape and if the lay-up tension is too high, delamination of the tape in its central zone has been observed in some cases by the Applicant, during the first centimeters of laying-up. In fact, the Applicant has observed that in such materials associating a tape of unidirectional fibers with each of its two faces carrying a thermoplastic veil, a preferential mechanical bond is established between the filaments located on the main faces of the tape and the veil, while the central zone of the tape, made up solely of filaments, corresponds to the zone with a lower shear strength.\nThat phenomenon may also be accentuated when a roller is used to deposit the tape. Under such circumstances, during the very first millimeters of bonding of the tape, the face in contact with the roller has a tendency to adhere to it, which further encourages delamination of the tape when its other face is then bonded to the preceding ply.\nThe Applicant has also observed the same phenomenon of shear of the yarn when laying up the tape along a trajectory that is curved in its plane, also known as “steering lay-up”. During this lay-up, the deposited tape undulates with an intensity that increases as the radius of lay-up decreases under the shear stress due to the surplus length of the filaments present on the internal radius of the tape compared with the filaments located on the outer edge of said tape."}
{"text": "Technical Field\nThe invention generally relates to compositions and methods for reprogramming non-pluripotent cells. More particularly, the invention relates to reprogramming mononuclear cells from low volume samples of blood.\nDescription of Related Art\nThrough their ability to self-renew and differentiate into any somatic cell type, induced pluripotent stem cells (iPSCs) have remarkable, but as yet untapped, therapeutic potential. Blood provides an attractive source of somatic cells for reprogramming as it is available in relatively large volumes and is easy to collect. However, blood, like other sources of explant-derived somatic cells for reprogramming, has some limitations in therapeutic applications since it must be HLA-typed and screened for pathogens before reprogrammed cells can be administered to a subject. What is needed in the art therefore is a readily accessible source of mononuclear cells for reprogramming that has been screened for HLA compatibility, pathogens and genetic abnormalities in the donor."}
{"text": "This invention relates to a solution of a substantially fully-imidized polyimidesiloxane in a substituted pyrrolidone solvent and to a method of coating substrates. In particular, it relates to solutions of polyimidesiloxanes in solvents containing N-cyclohexyl pyrrolidone (CHP), which do not develop a white crust when the solution is applied.\nA coating of a polymer can be made by dissolving the polymer in a solvent, spreading the solution over a substrate, and evaporating the solvent. When solutions of certain polymers in some solvents are spread and the solvent is evaporated, an objectionable white crust is formed. The whitening of the coating is a problem because it affects the flow characteristics of the solution in such a way that it sometimes does not uniformly coat the substrate.\nFor example, solutions of polyamic acids or polyimides in N-methylpyrrolidone (NMP) can develop this problem. The problem has been overcome for particular polyamic acids or polyimides by mixing the NMP with various co-solvents such as cyclohexanone (EP Application 364791 A2), tetramethylbenzene (EP Application 363737 A2), and butyl cellosolve acetate (Japanese Application 83-76689). While those solutions to the problem may be adequate for the particular polyimide being coated, they are not adequate for polyimidesiloxanes, because polyimidesiloxanes are less soluble than polyamic acids and may not dissolve in those mixtures of solvents. Also, whitening could still occur and there may be other reasons for not using those particular solvents. Until now a general solvent system for polyimidesiloxanes that does not produce whitening when a coating is formed has not been found."}
{"text": "This disclosure relates to data processing and data storage, and more specifically, to adaptation of one or more read voltage thresholds for a unit of data storage in a non-volatile memory system. Still more particularly, the disclosure relates a state-dependent read voltage threshold adaptation for a non-volatile memory system.\nNAND flash memory is an electrically programmable and erasable non-volatile memory technology that stores one or more bits of data per memory cell as a charge on the floating gate of a transistor or a similar charge trap structure. The amount of charge on the floating gate modulates the threshold voltage of the transistor. By applying a proper read voltage and measuring the amount of current, the programmed threshold voltage of the memory cell can be determined and thus the stored information can be detected. Memories storing one, two, three and four bits per cell are respectively referred to in the art as Single Level Cell (SLC), Multi-Level Cell (MLC), Three Level Cell (TLC), and Quad Level Cell (QLC) memories. In a typical implementation, a NAND flash memory array is organized in blocks (also referred to as “erase blocks”) of physical memory, each of which includes multiple physical pages each in turn containing a multiplicity of memory cells. By virtue of the arrangement of the word and bit lines utilized to access memory cells, flash memory arrays have generally been programmed on a page basis, but erased on a block basis.\nIn multi-level (i.e., MLC, TLC and QLC) NAND flash memory, information is stored by programming the memory cells to various quantized threshold voltage levels according to the device's programming algorithm, which maps the binary bit values to discrete threshold voltage levels. In response to a page read command, the binary bit values are retrieved by applying appropriate read voltages that divide the programmed threshold voltage window into discrete regimes and by then applying a reverse mapping between the detected threshold voltage levels and the corresponding binary bit values. Over the lifetime of a multi-level NAND flash memory device, the distributions of programmed threshold voltage generally become degraded due to the effects of wear on the memory cells. Consequently, it is generally desirable to adapt the read voltage thresholds defining the various bit values over time to compensate for the effects of wear and to extend the useful life of the NAND memory device."}
{"text": "FIELD OF THE INVENTION\nThe invention relates to a sheet-fed printing press, more particularly, having a plurality of stations arranged in a row and including at least one feeder, at least one digital printing unit and at least one delivery, a common sheet transport path, over which all sheets pass, extending between the at least one printing unit and the at least one delivery, and a postprocessing unit for the printing press, having a plurality of postprocessing stations arranged in a row.\nDigital printing units are understood to be printing units which can be activated pixel by pixel, such as laser printing units or inkjet printing units. Digital printing presses, which may contain a number of such printing units for multicolor printing, are able to print a new printed image for each new sheet, which permits the production of complete printed products, such as brochures, magazines or books, for example, at low outlay, even in the case of just a small number of copies. For this purpose, it is expedient and necessary, respectively, to be able to perform in one pass in the press all of the operations which follow the printing operation.\nPostprocessing stations for collating and further processing a block of sheets to form a brochure or the like have previously, in a logical development of conventional arrangements, been connected up rectilinearly to the normal processing line, so that the sheet transport direction, essentially during the entire processing operation, runs along a line which is more or less straight but is in any case disposed in a vertical plane. A printing press that is furnished or outfitted in this manner with a postprocessing unit for brochure printing cannot readily be converted for straightforward printing and vice versa, but rather, respective complex changeover and adjustment operations are necessary. In addition, the design of the printing press and that of the postprocessing unit depend upon one another to a great extent, for which reason it is generally not possible to utilize standard components."}
{"text": "1. Field of the Invention\nAn objective of the invention described herein relates to providing a device and method of securing metal coils to a flatbed trailer and desirably from ground level thereby obviating the need for anyone to be on the trailer itself at any time during the securing process.\n2. Description of Related Art\nIn the standard practice of securing metal coils to flatbed transport trailers, truck drivers and other persons often have to climb onto the flatbed to accomplish securing and tarping operations. This raises the possibility of injury as the truck driver and others may fall of the flatbed portion of the transport trailer and sustain injury and, in extreme circumstances, even death. Shippers have recognized the risk and safety concerns and in some cases have installed fall protection devices on the flatbed portion of the transport trailer often at an enormous expense. Some of known fall protection systems include mechanical devices that close in on the flatbed and have railings, fall prevention harnesses for the truck drivers so they can be hooked to a cable attached above the flatbed, and permanent flatbed extensions to allow for more footroom on the flatbed for the truck drive. While these device are improvements and in part address fall prevention issues, they tend to be expensive solutions and cumbersome for the user. More typically, fall prevention issues go completely unaddressed and simple mechanical load tie-down devices are found in the shipping art. Some these conventional mechanical devices are discussed immediately hereinafter.\nU.S. Pat. No. 2,212,180 to Murphy discloses a method of packaging heavy metal coils for shipment. The shipping arrangement disclosed by Murphy generally includes a sheet wrapping that fully encompasses a metal coil and simple metal bands are used to secure the sheet wrapping to the metal coil and an underlying support structure. Wooden planks are incorporated into the sheet wrapping for coil protection. The underlying support structure is used to lift the wrapped metal coil in conventional fashion.\nU.S. Pat. No. 3,154,026 to Klasing, Jr. discloses an articulated tie-down load control bracket which may be operatively associated with a metal coil in an “eye vertical” orientation. The control bracket includes two bracket pairs each comprising a pair of angles, a combined yoke and hook member common to both angles, and a pair of pivot pins attached to the combined yoke and hook member. The bracket pairs are located on opposite sides of the metal coil and are independently secured to a flatbed platform using chains and conventional brackets on the sides of the flatbed platform.\nU.S. Pat. No. 3,541,977 to Waldman discloses a pallet for transporting toroidal shaped articles comprising a base plate, a series of tie-downs that pass through the “eye” of the toroidal shaped article and secure the same to the base plate, and lifting cables that attach to the base plate for lifting the secured toroidal shaped article. U.S. Pat. No. 4,008,669 to Sumrell discloses an apparently similar mounting arrangement for mounting a metal coil to a pallet as that found in Waldman but discloses resiliently biased tie-down anchors for securing the metal coil and pallet to a flatbed platform.\nU.S. Pat. No. 4,204,479 to Rosa discloses a removable solid cover and containment arrangement for truck bed-mounted metal coils. The cover is rigid and generally semi-cylindrical in shape and is adapted to fit over the metal coil. The cover is further adapted so that after it has been placed over the metal coil a removable core member can be run through one side of the cover and the center of the metal coil and then firmly secured to another side of the cover so that any slippage of the metal coil will be effectively contained.\nU.S. Pat. No. 4,786,223 to Crissy et al. discloses a system for immobilizing vehicles having pneumatic tires during shipment utilizing tire chocks on each side of the tire and a flexible web harness passing over the tire upper circumference having ends affixed to the chocks. A winch mounted upon one chock tensions the harness while a quick-release fitting at the other end of the harness attaches to the other chock.\nU.S. Pat. No. 5,888,039 to Cooley discloses cargo securing system including several elongated block portions having a plurality of studs extending downwardly therefrom for coupling with through holes in a flatbed trailer. A plurality of pipe attachment portions are also provided with each having a recess extending within a lower surface thereof. The recess is dimensioned for coupling with the block portion. A pipe harness is provided comprising a pair of elongated straps, a pair of short cross straps, and crisscrossing central straps. The short cross straps have free ends with buckles disposed thereon. An extension strap is adapted for coupling with free ends of the elongated straps of the pipe harness. The extension strap has a clip on a free end thereof for engaging existing securement buckles on the trailer. A coupling strap is provided having clips disposed on opposing free ends thereof. The coupling strap couples with the buckles of the short cross straps and the existing securement buckles of the flatbed trailer.\nU.S. Pat. No. 7,214,014 to Stanley discloses a cargo restraint system comprising a tie-down assembly of restraining tie-down arms having interconnected common proximal extremities and a plurality of distal extremities. The tie-down assembly is arranged such that an angle is formed between each arm, the vertex of the angle being congruent with the common proximal extremities.\nU.S. Pat. No. 7,270,507 to Jernigan discloses a safety coil sling for securing a metal coil on a flatbed trailer having a round or elliptical primary cable attached at one end to a first anchoring chain and at an opposite end to a second anchoring chain, and a plurality of connecting cables for connecting opposite sides of the primary cable across the top of the coil, thereby securing the coil on the trailer. Another disclosed embodiment includes the primary cable being divided into parallel segments at the points of attachment of the connecting cables and the connecting cables are slidably attached to one of the parallel segments to allow the connecting cables to slide within a fixed area, thus providing a mechanism for relieving tension on the sling during weight shifts within the coil occurring during transport.\nWhile numerous load securing device are known in the shipping art as evidenced by the foregoing patents, there is generally a lack of attention to fall prevention issues and further, ease of use, or lack thereof, of such devices in practice. The focus of the prior art in the shipping field is generally on securing the load with little regard given to the safety of the user and labor required to effect installation of the load securing device on an intended load. This disclosure now turns to devices and methods for efficient, labor-saving, and safe load securing operations involving metal coils."}
{"text": "Transparent display technology is an emerging display technology. When a transparent display device display an image, objects disposed behind the transparent display device also can be seen through the device. That needs to enhance a transmittance of the transparent display device without losing its color gamut and contrast.\nMeanwhile, the transparent display technology is a typical representative of the green display technology. For example, a conventional liquid crystal display needs a backlight module as a light source to display an image. The transparent display may use an ambient light to meet the need of the backlight. Especially in the daytime, the transparent display almost needs no backlight. The power consumption of the transparent display is about one-tenth of the ordinary liquid crystal display.\nIn order to improve the transmittance of the transparent display device, the original three primary colors of RGB color filter is usually reduced to 1/3 colors. However, by doing so, the color saturation of the image is also decreased.\nTherefore, it needs to develop a transparent display device with high color saturation and high transmittance.\nIn order to increase the brightness of an image displayed by a transparent display device, an optical sheet can be placed under a light guide plate thereof, to reflect the light emitted from the light guide plate to the bottom and reflect it back to the light guide plate. The light guide plate includes a first refractive layer and a second refractive layer to allow a natural light transmits through the light guide plate. The first refractive layer is located adjacent to the liquid crystal panel. The second refractive layer has a higher refractive index than the first refractive layer. A first polarized light in the light guide plate is total reflected by a interface between the first refractive layer and the second refractive layer, and incident on the liquid crystal panel above the light guide plate, and a natural light transmits through the first and the second refraction layers incident on the liquid crystal panel, so that the user in front of the display device can see an object behind the light guide plate. However, the addition optical sheet and the special design of the light guide plate will increase the manufacturing cost of the display device."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a printed circuit board. More particularly, the present invention relates to a printed circuit board integrated with a two-axis fluxgate sensor and a method for manufacturing the same.\n2. Description of the Related Art\nThe existence of magnetic energy has been proven through various physical phenomena. For example, a fluxgate sensor enables a human to indirectly perceive magnetic energy, since magnetic energy is not perceivable by human sense organs, such as eyes or ears. Regarding a fluxgate sensor, a magnetic sensor employing a soft magnetic coil has been conventionally used. A magnetic sensor is made by winding a coil around a relatively large bar-shaped core or annular-shaped core, which is formed of a soft magnetic ribbon. In addition, an electronic circuit is employed to obtain a magnetic field in proportion to the measured magnetic field.\nThe conventional fluxgate sensor, however, has the following problems, some of which are due to the structure of the conventional fluxgate sensor. In the structure of a conventional fluxgate sensor the coil being wound around the large bar-shaped core or annular-shaped core, which is made of the soft magnetic ribbon, results in high production costs and a large volume of the overall system.\nIn addition, flux leakage is inevitable in the flux change due to the excitation coil and the detected magnetic field. Accordingly, high sensitivity is not readily achieved.\nFurther, a conventional fluxgate sensor for detecting magnetic energy in two axes includes two separate one-axis fluxgate sensors arranged in perpendicular relation with each other, which results in yet another problem of the device size increasing."}
{"text": "Storage library systems are often used by enterprises and the like to efficiently store and retrieve data from storage media. In the case of some storage libraries, the media are data cartridges (e.g., tape cartridges) that are typically stored and indexed within a set of magazines. When particular data is requested, a specialized robotic mechanism finds the appropriate cartridge, removes the cartridge from its magazine, and carries the cartridge to a drive that is designed to receive the cartridge and read its contents. Some storage libraries have multiple drives that can operate concurrently to perform input/output (IO) operations on multiple cartridges.\nTo operate properly, the robotic mechanisms are expected to reliably (e.g., repeatably and accurately) and rapidly find, retrieve, and deliver desired cartridges throughout the storage library cartridge inventory. This functionality can be facilitated by configuring the robotic mechanism to move a hand assembly in at least three axes (e.g., x, y, and z directions, and sometimes one or more of pitch, roll, or yaw), and to include one or more sensors to reliably detect the position and/or orientation of the hand assembly. When in its desired location, the hand assembly is activated to reliably grip the desired cartridge and remove it from a magazine or drive, or to reliably release the cartridge into a desired magazine slot or drive. The gripping and/or releasing of the cartridge is typically performed by a gripper that is part of the hand assembly.\nMany traditional implementations of the gripper include complex, expensive, and/or heavy components. For example, some designs use solenoid-operated opposing gripper plates acting against a spring load. While these and other techniques can be reliable, they tend to manifest undesirable power consumption characteristics and to include heavy and/or expensive components (e.g., solenoids, power cables, specialized hardware and software, etc.)."}
{"text": "The invention relates to a switch device which includes a cam groove and a lock pin, and has the function of holding an operating member in two positions.\nA so-called push lock-type switch device has the following construction. A cam groove, having, for example, a heart-shape, is formed in one of a device body and an operating member movably mounted on the device body, while a lock pin is provided at the other, and the lock pin slides in the cam groove in accordance with the movement of the operating member so as to hold the operating member in two positions (an original position and a pushed-in position). A holding spring for urging the lock pin toward the cam groove is provided together with this lock pin.\nFIG. 6 shows the portion of the operating member of the above conventional construction including the lock pin and the holding spring. The operating member 51 is mounted on the device body (not shown) for reciprocal movement in directions of arrow Q. A proximal end portion 52a of the lock pin 52 is fitted in a fitting hole 51a formed in the operating member 51. A distal end portion 52b of the lock pin 52 is held in sliding contact with the cam groove (not shown). The holding spring 53 has a linear shape, and has opposite end portions fitted respectively in fitting grooves 51b and 51b. The holding spring 53 serves to urge the lock pin 52 toward the cam groove.\nIn this construction, when mounting the operating member 51 on the device body (not shown), there is a fear that the lock pin 52 drops, and the assembling operation must be carried out while taking care not to allow the lock pin to drop. And besides, since this switch device itself is small, there has been encountered a disadvantage that the efficiency of the assembling operation is low.\nThis invention has been made under the above circumstances, and an object of the invention is to provide a switch device in which a device body and an operating member can be assembled together while preventing the dropping of a lock pin, and besides the lock pin can be moved in a predetermined manner after the assemblage.\nIn order to solve the aforesaid object, the invention is characterized by having the following arrangement.\n(1) A switch device comprising:\na device body;\nan operating member, for operating a switch portion by movement thereof, reciprocally movably mounted on the device body;\na cam groove portion formed in one of the device body and the operating member;\na lock pin which includes a proximal end portion mounted on the other of the device body and the operating member, and a distal end portion which can slide in the cam groove portion so as to releasably hold the operating member in two positions;\na holding spring for urging the lock pin toward a bottom surface of the cam groove portion, mounted on the other of the device body and the operating member with its opposite end portions fitted therein;\na pin engagement portion, for engaging the lock pin to prevent the dropping of the lock pin in a condition before the operating member is mounted, formed on the other of the device body and the operating member; and\nan engagement cancellation portion, for disengaging the lock pin from the pin engagement portion in a mounted condition of the operating member, formed on the one of the device body and the operating member.\n(2) The switch device according to (1), wherein a holding spring engagement portion for preventing the dropping of the holding spring in the condition before mounting the operating member is formed on the other of the device body and the operating member.\n(3) The switch device according to (1), wherein\nthe lock pin is adapted to rotate about the proximal end portion,\nthe distal end portion of the lock pin is moved within a movement range of the cam groove when the operating member performs the switching operation, and\nthe pin engagement portion is provided outside the movement range.\n(4) The switch device according to (1), wherein\nan engagement cancellation portion for disengaging the lock pin from the pin engagement portion is formed on the one of the device body and the operating member.\n(5) The switch device according to (4), wherein\nwhen the operating member is assembled to the device body, the engagement cancellation portion disengages the lock pin from the pin engagement portion and move the lock pin to the movement range."}
{"text": "1. Field\nExemplary embodiments of the present invention relate to a semiconductor device fabrication technology, and more particularly, to an image sensor, and a method for fabricating the image sensor.\n2. Description of the Related Art\nAn image sensor is a device that transforms optical images of one or more dimensions on a photographic subject into electric signals. An image sensor includes a pixel array region for receiving incident light and generating photocharges and a logic region for generating electric signals by processing the generated photocharges. The logic region may also be called a peripheral circuit region. Incident light radiates not only upon the pixel array region but also upon the logic region as well. The photocharges generated in the logic region by the incident light radiating thereon may act as noise on the electric signals generated in the logic region, that is, electric signals generated by processing the photocharges generated in the pixel array region. Therefore, the logic region may use a light shielding pattern for blocking out the incident light.\nGenerally, a light shielding pattern is formed over a substrate of the logic region. Therefore, there is a step height difference between the pixel array region and the logic region. The step height difference between the two regions that is caused by the light shielding pattern deteriorates the characteristics and integration of the image sensor."}
{"text": "1. Field of the Invention\nThe present invention relates to compressible supports for mine roofs and the like. More specifically, the present invention relates to vertical support columns having compressible supports.\n2. Brief Description of the Prior Art\nSupport columns are used to counteract the force that gravity exerts on the mass of an object. For example, vertical support columns are positioned between two horizontal objects to keep the objects spaced away from one another.\nIn underground mining, vertical support columns are used as primary or secondary supports for mine roofs. Due to the tremendous forces exerted on the vertical support columns by the earth and rock above the mine roof, including forces resulting from shifting and settling of the overhead earth and rock, rigid vertical support columns can bend or break. Complete failure of a vertical support column can cause an isolated collapse of the mine roof. Therefore, compressible support columns, which yield during settling or shifting of the mine roof, are often used.\nIn mining operations, compressible vertical support columns are generally positioned perpendicularly between a mine roof and a mine floor. One type of compressible support column is disclosed in U.S. Pat. No. 4,052,029 to Townsend. The Townsend patent uses telescoping members and compressible materials to reduce the overall length of the support column as the earth and rock above the mine roof settles or shifts. Specifically, the Townsend patent discloses a support column having a hollow lower steel member and a hollow upper steel member, where the upper steel member fits over and slideably telescopes toward the lower steel member. The hollow cavity within the lower steel member is completely filled with wood positioned with the grain of the wood oriented parallel to a vertical axis of the assembled column. The hollow cavity in the upper steel member is only partially filled with wood, with the grain of the wood also aligned with the vertical axis of the assembled column. The space left between the wood in the upper and lower steel members is filled with discs, as needed, to adjust the length of the assembled compressible support column between a mine roof and a mine floor.\nWhen earth and rock above the compressible support column disclosed in the Townsend patent shift, settle, or are otherwise subjected to additional force, the wood inside of the upper and lower steel members compresses. The compression reduces the overall length of the assembled compressible support column, easing the force being exerted on the column. Additional settling or shifting of the mine roof further compresses the load resisting material, until the lower steel member is telescoped to its full extent inside of the upper steel member. At this point, the upper and lower steel columns begin to compress, causing the upper and lower steel members to bulge outwardly, away from the vertical axis of the assembled upper and lower steel columns.\nOne disadvantage of the prior art compressible support columns is that they are expensive to make. Another disadvantage is that adjusting the length of the columns during installation is time consuming. Therefore, one object of the present invention is to provide a column having a compressible support that is inexpensive to manufacture. Another object of the present invention is to make a column having a compressible support that can be installed quickly.\nThe present invention generally includes a compressible support which includes a male member defining a plurality of ridges and grooves and a collar which has an internal surface defining a plurality of corrugations. The corrugations adjustably engaging grooves defined by the male member, the ridges defined by the male member, and the corrugations defined by the collar yield when the male member and collar are compressed together\nOne embodiment of the present invention generally includes a post and a collar. A cap and a base may also be provided. The post has a first end, a second end, and a longitudinal axis. The post also includes ridges and grooves, generally in the form of threads, positioned adjacent the first end of the post.\nThe collar has a first end, a second end, an internal surface, and an external surface. The second end of the collar is positioned adjacent the first end of the post during installation. Corrugations made from a material harder than the ridges and grooves of the post are positioned adjacent the internal surface of the collar. The corrugations adjustably engaging the ridges and grooves of the post, which are generally in the form of threads, are adjacent the first end of the post. The corrugations and threads allow the collar to be adjustable along a longitudinal axis of the post, essentially by threading the collar onto or off of the first end of the post. An optional cap and an optional base may be positioned adjacent the ends of the post and collar to distribute the force applied to the post and collar over a greater surface area.\nWhen an initial or preloading force is exerted on the male member or post and the collar, referred to hereafter as the compressible support, the compressible support provides an equal and opposite force. As the force on the compressible support increases, the ridges positioned adjacent to the first end of the male member crush, yield, fracture, or strip, allowing the collar to move toward the second end of the male member, decreasing the overall length of the compressible support. Additional force causes additional crushing of the grooves adjacent the first end of the male member, until the movement of the collar toward the second end of the male member is arrested. At this point, further force causes the first end of the male member and the collar to compress, further decreasing the overall length of the compressible support.\nAnother embodiment of a compressible support according to the present invention generally includes a male member, a collar, and a post. In this embodiment, the male member is separate from the post.\nThese and other advantages of the present invention will be clarified in the description of the preferred embodiments taken together with the attached drawings in which like reference numerals represent like elements throughout."}
{"text": "The detection of Cryptosporidium oocysts in water presently relies on the concentration of particulate matter from large volumes of water prior to staining with fluorescently labelled monoclonal antibodies. Until recently, detection and identification of fluorescently labelled oocysts required examination of the sample using epifluorescence microscopy. The tedious and labour intensive nature of this detection method,. in particular the amount of fluorescent microscopy required, limited the monitoring work which could be performed. The development of flow cytometric assisted detection methods has alleviated some of these problems and enabled the routine monitoring of water for the presence of Cryptosporidium oocysts (Vesey et al., 1994A). However, a major limitation of all these methodologies is the lack of oocyst viability measurements. Methods to determine viability such as animal infectivity and excystation, are impractical because of the low number of oocysts normally present in water samples and the tedious nature of such tests.\nThe presence of dead Cryptosporidium oocysts in drinking water is of little significance to public health, however if oocysts are viable the risk to public health is enormous. Moreover, of the seven species of Cryptosporidium; C. parvum, C. muris, C. meleagridis, C. serpentis, C. nasorum, C. wrairi and C. baileyi only one, C. parvum is infectious to humans (Rose et al., 1988), and yet all the currently available monoclonal antibodies are not species specific. There is therefore, an urgent requirement to develop an effective method for determining the viability of C. parvum oocysts in water.\nFluorescence in situ hybridisation (FISH) is a relatively new method by which microorganisms can be specifically labelled. The technique is reliant upon the identification of a specific sequence of nucleic acid within the target organism. Probes targeting a specific nucleic acid sequence are then synthesised and labelled with a fluorochrome. The cell is then permeabilised and the complementary sequence allowed to hybridise with the target sequence resulting in specific labelling of the target cell.\nThe use of ribosomal RNA (rRNA) targeted oligonucleotide probes with FISH and flow cytometry has been reported by Amann et al. (1990B). Molecules of rRNA are ideal targets for fluorescently labelled nucleic probes for several reasons: (1) high sensitivity can be achieved since the target molecules are present in very high numbers; (2) a denaturation step is not required during the procedure as the target region is single stranded; and (3) rRNA has a short half life and will only be present in a high copy number in viable cells."}
{"text": "The present invention is directed generally to catalytic antibodies and, more particularly, to a novel method of producing same. The method of the present invention is predicated in part on the exploitation of the products of catalysis to induce B cell mitogenesis. In a preferred embodiment, a growth factor having an ability to induce B cell mitogenesis is linked to a target antigen to which catalytic antibodies are sought. B cell mitogenesis is then dependent on the catalytic cleavage of the antigen portion of the growth factor by catalytic antibodies on the surface of B cells. The method of the present invention is useful for generating catalytic antibodies for both therapeutic and diagnostic purposes.\nSequence Identity Numbers (SEQ ID NOs.) for the nucleotide and amino acid sequences referred to in the specification are defined following the Examples.\nThroughout this specification, unless the context requires otherwise, the word xe2x80x9ccomprisesxe2x80x9d, or variations such as xe2x80x9ccomprisesxe2x80x9d or xe2x80x9ccomprisingxe2x80x9d, will be understood to imply the inclusion of a stated element or integer or group of elements or integers but not the exclusion of any other element or integer or group of elements or integers.\nThe rapidly increasing sophistication of recombinant DNA technology is greatly facilitating research and development in the medical and allied health fields. A particularly important area of research is the use of recombinant antigens to stimulate immune response mechanisms and outcomes. However, until now, recombinant techniques have not been particularly effective in the generation of catalytic antibodies.\nCatalytic antibodies are highly substrate specific catalysts which can be used, for example, to proteolytically activate or inactivate proteins. Catalytic antibodies have great potential a therapeutic agents in human diseases such as rheumatoid arthritis, AIDS and Alzhiemer\"\"s disease amongst many others.\nAntibody therapy is already routinely used in patients. Antibodies have a half-life of about 23 days in the circulation of humans which is a clear advantage over other drugs. Catalytic antibodies, however, are considered to be even more effective. They are recycled after their antigenic encounter and are not bound to the antigen as occurs with xe2x80x9cclassicalxe2x80x9d antibodies. Catalytic antibodies should, therefore, function at a much lower dose than classical antibodies and could be used at sub-immunogenic doses. Catalytic antibodies would be particularly useful in long term therapy.\nTraditionally, catalytic antibodies have been generated by immunising mice with transition state analogs. Such antibodies have been shown to catalyse several chemical reactions. However, this approach has a severe limitation in that it is difficult to predict the structure of transition state analogs which effect proteolysis of specific proteins. Immunising a mouse with a transition state analog is by definition inefficient since it selects B cells on the ability of surface immunoglobulins to bind the analogs and not on the catalytic activity of the surface immunoglobulins. This is one of the reasons why catalytic antibodies have relatively low turn-over rates and cannot compete with the naturally occurring enzyme counterparts. As a consequence, catalytic antibodies have not previously achieved prominence as therapeutic or diagnostic tools.\nThere is a need, therefore, to develop a more efficacious approach to generating catalytic antibodies having desired catalytic specificity. In accordance with the present invention, the inventors have developed such an approach based on a recombinant or a synthetic growth factor having an ability to induce B cell mitogenesis. A precursor form of the growth factor selects xe2x80x9ccatalyticxe2x80x9d B cells. The present invention provides, therefore, for the exploitation of the products of catalysis for B cell activation which may and can be antigen binding site independent.\nAccordingly, one aspect of the present invention is directed to a recombinant or synthetic growth factor or a precursor thereof comprising a B cell surface molecule binding portion wherein said growth factor or a catalytic product of said precursor is capable of inducing B cell mitogenesis.\nMore particularly, the present invention provides a recombinant or synthetic growth factor comprising a B cell surface molecule binding portion wherein said growth factor induces B cell mitogenesis.\nIn one aspect of the present invention, the recombinant or synthetic growth factor comprises a B cell surface molecule binding portion and a portion providing T cell dependent help for a B cell such that said growth factor induces B cell mitogenesis.\nIn another aspect, the portion providing T cell dependent help for a B cell is supplied independently of the growth factor. An example of an exogenously supplied portion having T cell dependent help for a B cell is anti-CD40 antibodies or functional equivalents thereof.\nIn a particularly preferred embodiment, the B cell surface molecule binding portion comprises a B cell surface immunoglobulin binding portion although the present invention extends to a range of B cell surface molecules the binding, interaction and/or cross-linking thereof leads to or facilitates B cell mitogenesis. Reference hereinafter to a xe2x80x9cB cell surface moleculexe2x80x9d includes reference to a B cell surface immunoglobulin. The portion providing T cell dependent help for a B cell is preferably but not exclusively a T cell epitope. Reference hereinafter to a portion providing T cell dependent help for a B cell includes a T cell epitope.\nThe present invention further contemplates a composition of matter capable of inducing B cell mitogenesis said composition of matter comprising components selected from:\n(i) a recombinant or synthetic molecule comprising a B cell surface molecule binding portion;\n(ii) a recombinant or synthetic molecule in multimeric form comprising a B cell surface molecule binding surface molecule binding portion;\n(iii) a recombinant or synthetic molecule of (i) or comprising a further portion providing T cell dependent help for a B cell; and\n(iv) separate compositions mixed prior to use or used sequentially or simultaneously comprising in a first composition a component having a B cell surface molecule binding portion and in a second composition a molecule capable of providing T cell dependent help for a B cell.\nPreferably, the molecule capable of providing T cell dependent help for a B cell is a T cell epitope or is an anti-CD40 antibody or functional equivalents thereof.\nPreferably, to ensure cross-liking of B cell surface molecules to induce blastogenesis, the growth factor comprises at least two B cell surface molecule binding portions. Alternatively, where the growth factor is present in multimeric form, the molecule may comprise a single B cell surface molecule binding portion.\nEven more preferably, the growth factor comprises a T cell epitope portion flanked by, adjacent to or proximal with at least one B cell surface molecule-binding portion. The presentation of a T cell epitope on the surface of a B cell allows for B cell mitogenesis. The term xe2x80x9cB cell mitogenesisxe2x80x9d is used herein in its broadest context and includes B cell activation, clonal expansion, affinity maturation and/or antibody secretion as well as growth and differentiation. The term xe2x80x9cmitogenicxe2x80x9d as used herein means xe2x80x9cmitogenesisxe2x80x9d.\nIn another embodiment, the recombinant or synthetic growth factor comprises a further portion permitting multimerisation of the growth factor. A multimer comprises two or more growth factor molecules or a precursor thereof. Examples of portions inducing multimerisation include but are not limited to an antibody, a region facilitating formation of cross-linked molecules or a signal peptide. Cross-linkage in this context includes any interaction that provides bonding adequate to lead to multimer formation including but not limited to covalent linkage, ionic linkage, lattice association, ionic bridges, salt bridges and non-specific molecular association. A particularly preferred embodiment of the present invention is directed to the use of a signal peptide such as the signal peptide of the Outer Membrane Portion A (ompA) [Skerra, Gene, 151: 131-135, 1994] or a functional derivative thereof A xe2x80x9cfunctional derivativexe2x80x9d in this context is a mutant or derivative of the ompA signal peptide (or its functional equivalent) which still permits multimer formation of the growth factor.\nIn a preferred embodiment, where the B cell surface molecule is an immunoglobulin, the B cell surface binding portion of the growth factor generally binds to part of an immunoglobulin such as the variable portion of a heavy or light chain of an immunoglobulin.\nAn example of an immunoglobulin binding molecule is protein L from Peptostreptococcus magnus. Protein L has five immunoglobulin-binding domains. Other immunoglobulin binding molecules include protein A, protein G and protein H. The present invention, however, extends to any molecule capable of binding to a B cell surface component including, for example, a ligand of a B cell receptor.\nThe portion of the recombinant or synthetic molecule defining a T cell epitope is presented to a preferably already primed T cell to induce B cell proliferation and affinity maturation of an antibody in the germinal centre. This is generally accompanied by immunoglobulin class switching and antibody secretion into the serum. Generally, the T cell epitope is internalised within the B cell and presented on major histocompatbility complex (MHC) class II.\nIn a particularly preferred embodiment, the T cell epitope is from hen egg lysozyme (HEL) [Altavia et al, Molecular Immunology, 31: 1-19, 1994] or is a derivative thereof such as a peptide comprising no acids 42 to 62 from HEL or a homologue or analog thereof. This T cell epitope is recognised by the T cell receptor (TCR) of HEL specific T cells when presented by an antigen presenting cell (APC) on the MHC class II molecule H-2Axcexa in mice or other MHC class II molecules or their equivalents in other mammals such as humans. Examples of other T cell epitopes include but are not limited to tetanus toxoid, ovalbumin, malarial antigens as well as other regions of HEL. One skilled in the art would readily be able to select an appropriate T cell epitope.\nIn an alternative embodiment, the portion providing T cell dependent help of a B cell is not a T cell epitope in the classical sense but is nevertheless able to function in a similar manner. An example of such a portion is an anti-CD40 antibody.\nAnother aspect of the present invention contemplates a recombinant or synthetic molecule having the structure:\nX1X2X3\nwherein:\nX1 and X3 may be the same or different and each is a B cell surface molecule binding entity;\nX2 is a portion providing T cell dependent help for a B cell; and wherein the recombinant or synthetic molecule is capable of inducing T cell dependent B cell mitogenesis of the B cell which X1 and X3 bind.\nThe representation X1X2X3 is not to be taken as imposing any sequential constraints and the present invention covers any sequence comprising the elements X1X2X3 such as X1X3X2, X3X2X1, X2X3X1, X2X1X3, and X3X1X2.\nThe present invention further provides a recombinant or synthetic molecule having the structure:\nX1X2X3\nwherein:\nX1 and X3 may be the same or different and each is a B cell surface immunoglobulin binding entity;\nX2 is a T cell epitope; and\nwherein the recombinant or synthetic molecule is capable of inducing T cell dependent B cell mitogenesis of the B cell to which X1 and X3 bind.\nThe present invention yet further provides a recombinant or synthetic molecule having the structure:\nxe2x80x83MX1[X2]b\nwherein:\nX1 is a B cell surface molecule binding entity such as a B cell surface immunoglobulin binding entity;\nX2 is a portion providing T cell dependent help for a B cell;\nM is a portion enabling or facilitating multimer formation of the recombinant or synthetic molecule; and\nb represents the number of X2 molecules and may be 0, 1 or  greater than 1.\nAn example of M is a signal peptide such as from ompA. The multimeric effect of M is to provide a molecule [MX1[X2]b]k where M, X1, X2 and b are as defined above and k is 1 or  greater than 1. An example of k is from about 2 to about 1000 and more preferably from above  greater than 1 to about 50-200.\nThe optional nature of X2 may alternatively relate to X1. The molecular components of MX1[X2]b may be in any order.\nPreferably, but not exclusively X1 is derived from the L protein from P. magnus. \nPreferably, X2 is a T cell epitope derived from HEL such as but not limited to amino acids 42 to 62 from HEL. In an alternative embodiment, where b=0 X2 is not a T cell epitope and is exogenously supplied facilitating T cell dependent help for a B cell. In this case, X2 may be, for example, an anti-CD40 antibody.\nIn a particularly preferred embodiment, X1X2X3 as defined above is referred to herein as xe2x80x9cLHLxe2x80x9d where L is an immunoglobulin binding entity from L protein from P. magnus and H corresponds to amino acids 42 to 62 from HEL or an equivalent thereof. The nucleotide and amino acid sequences of LHL are shown in SEQ ID NO:1 and SEQ ID NO:2, respectively.\nThe present invention extends to a growth factor or precursor thereof having the amino acid sequence set forth in SEQ ID NO:2 or an amino acid sequence having at least 60% similarity thereto (excluding the amino acid sequence of an antigen) or an amino acid sequence encoded by nucleotide sequence set forth in SEQ ID NO:1 or a nucleotide sequence capable of hybridizing to SEQ ID NO:1 under low stringency conditions.\nReference herein to a low stringency at 42xc2x0 C. includes and encompasses from at least about 1% v/v to at least about 15% v/v formanide and from at least about 1M to at least about 2M salt for hybridisation, and at least about 1M to at least about 2M salt for washing conditions, Alternative stringency conditions may be applied where necessary, such as medium stringency, which includes and encompasses from at least about 16% v/v to at least about 30% v/v formamide and from at least about 0.5M to at least about 0.9M salt for hybridisation, and at least about 0.5M to at least about 0.9M salt for washing conditions, or high stringency, which includes and encompasses from at least about 31% v/v to at least about 50% v/v formamide and from at least about 0.01M to at leas. about 0.15M salt for hybridisation, and at least about 0.01M to at least about 0.15M salt for washing conditions.\nThe two xe2x80x9cLxe2x80x9d domains from protein L linked by the xe2x80x9cHxe2x80x9d domain from HEL build the growth factor LHL which crosslinks the surface immunoglobulin on B cells. This cross-linking induces B cell activation and blast formation, The internalisation and processing of LHL leads to the presentation of H on MHC II. T cell recognition of MHC II with the H peptide signals the activated B cell to proliferate and undergo antibody class switching and secretion.\nThe mitogenic growth factor of the present invention is most useful in generating antibodies of desired catalytic specificity when in a precursor form which selects xe2x80x9ccatalyticxe2x80x9d B cells. A precursor growth factor comprises a target antigen to which a catalytic antibody is sought and which mask antigen-independent clonal expansion of B cells. Upon cleavage of the antigen by a selected B cell surface immunoglobulin, the growth factor can induce B cell mitogenesis.\nIn effect, then B cells are selected on the catalytic activity of their surface immunoglobulin rather than on their binding to a transition state analog. This allows for affinity maturation in the germinal centres and ensures xe2x80x9ccatalytic-maturationxe2x80x9d to obtain the highest enzymatic turn-over rate possible in vivo. This aspect of the present invention is achieved by designing a B cell growth factor precursor comprising X1X2X3 as defined above, such as LHL, shielded and substantially inactive until released through cleavage by a catalytic antibody on a B cell surface. The term xe2x80x9ccleavagexe2x80x9d in this context is not limiting to the breaking of bonds but includes interaction adequate to remove or reduce shielding of the B cell growth factor.\nThe liberated X1X2X3activates the catalytic B cell by crosslinking the B cell surface molecules such B cell surface immunoglobulins via the X1X3 domains. X1X2X3 is then internalised and processed analogous to a normal antigen. Intracellular processing generates the T cell epitope X2 incorporated in X1X2X3. The presentation of X2 on the B cell surface leads to T cell dependent clonal expansion of the B cell as well as catalytic maturation and secretion of the catalytic antibody. The catalytic antibodies can then be detected in serum and catalytic B cells can be recovered by standard techniques.\nAccording to this aspect of the present invention, there is provided a recombinant or synthetic growth factor precursor comprising an antigen linked or otherwise associated with a growth factor capable of inducing B cell mitogenesis, wherein the products of catalysis of said growth factor precursor permit B cell mitogenesis which may and can be antigen binding site independent.\nAccording to another aspect of the present invention, there is provided a recombinant or synthetic growth factor precursor comprising an antigen linked to or otherwise associated with a B cell surface binding molecule portion and a portion providing T cell dependent help to a B cell such that upon cleavage of said antigen or of a region proximal to said antigen, said B cell surface molecule portion and said T cell dependent help portion form a growth factor which induces B cell mitogenesis.\nPreferably, the present invention is directed to a recombinant or synthetic growth factor precursor comprising an antigen linked to or otherwise associated with a B cell surface immunoglobulin binding portion and a T cell epitope portion such that upon cleavage of said antigen or of a region proximal to said antigen, said B cell surface immunoglobulin and T cell epitope form a growth factor having B cell mitogenic properties.\nIn an alternative embodiment, the present invention provides a recombinant or synthetic growth factor precursor comprising an antigen linked to or otherwise associated with a B cell surface molecule binding portion and a portion conferring a multimeric form of said growth factor precursor when cleavage of said antigen or of a region proximal to said antigen provides a growth factor comprising a B cell surface molecule binding portion having the ability to induce B cell mitogenesis.\nThe antigen according to this aspect of the present invention is any antigen to which a catalytic antibody is sought. Examples include cytokines such as but not limited to tumor necrosis factor (TNF), an interleukin (IL) such as IL-1 to IL-15, interferons (IFN) such as IFNxcex1, IFNxcex2 or IFNxcex3, colony-stimulating factors (CSF) such as granulocyte colony-stimulating factor (G-CSF), granulocyte-macrophase colony-stimulation factor (GM-CSF), blood factors such as Factor VIII, erythropoietin and haemopoietin, cancer antigens, docking receptors from pathogenic viruses such as HIV, influenza virus or a hepatitis virus (eg. HEP A, HEP B, HEP C or HEP E) and amyloid plaques such as in Alzheimer\"\"s disease patients or myeloma patients. More particularly, in the case of TNF, proteolytic inactivation of TNF would be useful in the treatment of rheumatoid arthritis and toxic shock syndrome. By targeting viral docking receptors, pathogenic viruses such as HIV, hepatitis viruses and influenza viruses are rendered effectively inactive. Catalytic antibodies will also be useful in the clearance of amyloid plaques in Alheimer\"\"s disease or myeloma disease patients. Targeting IgE, for example, may provide a mechanism for treating inflammatory conditions such as asthma.\nThe catalytic antibodies of the present invention may also be useful in detoxifying drugs such as drugs consumed by an individual. For example, the effects of cannabis or heroin or other drugs could be treated in an individual by the administration of catalytic antibodies directed to the active components of those drugs. Furthermore, catalytic antibodies may be useful in the treatment of autoimmune and inflammatory disease conditions such as by targeting autoimmune antibodies. Catalytic antibodies also have a use in environmental situations and could be directed to environmental pollutants such as petroleum products and plastics. In all these situations, suitable antigens would be selected and incorporated into the growth factor precursor of the present invention.\nIn a related aspect of the present invention, the xe2x80x9cantigenxe2x80x9d portion of the growth factor precursor can be mimiced by a target site such as an amino acid linker sequence comprising a protease cleavage site. Examples include an amino acid linker sequence comprising the tabacco etch virus (TEV) protease cleavage site More particularly, in the case of a TEV protease cleavage site, cleaving of the amino acid linker sequence by the TEV protease would be useful for producing characteristics similar to those of a catalytic antibody. This provides a useful model system for developing growth factor molecules.\nAnother aspect of the present invention provides a recombinant or synthetic molecule having the structure:\nAX1X2X3A\nwherein\nA is a target antigen for which a catalytic antibody is sought;\nX1 and X3 may be the same or different and each is a B cell surface molecule binding entity;\nX2 is a portion providing T cell dependent help for a catalytic B cell;\nwherein a catalytic antibody on the surface of said B cell is capable of cleaving all or part of A from said recombinant or synthetic molecule resulting in the molecule [Axe2x80x2]X1X2X3[Axe2x80x2] wherein Axe2x80x2 is optionally present and is a portion of A after cleavage with the catalytic antibody and wherein said resulting molecule is capable of inducing T cell dependent clonal expansion of the B cell to which X1 and X3 bind.\nThe AX1X2X3A molecule may be in any sequence order.\nAnther aspect of the present invention is directed to a recombitant or synthetic molecule having the structure:\n[M]cAX1[X2]d[X3]eA[M]f\nwherein:\nA is a target antigen for which a catalytic antibody is sought;\nX1 and X3 may be the same or different and each is a B cell surface molecule binding entity;\nX2 is a portion providing T cell dependent help for a B cell;\nM is a portion enabling or facilitating multimer formation of the recombinant or synthetic molecule;\nc and f may be the same or different and each is 0, 1 or  greater than 1;\ne is 0 or 1 with the proviso then if both c and f are 0 then e cannot be 0;\nd is 0 or 1 or  greater than 1;\nwherein a catalytic antibody on the surface of said B cell is capable of cleaving all or a portion of A from said recombinant or synthetic molecule and a resulting catalytic product is capable of inducing B cell mitogenesis.\nThe growth factor precursor enables an antigen to be recognised by a B cell via a growth factor capable of inducing B cell mitogenesis. The antigen effectively masks B cell activation without prior catalytic activation. The growth factor is in xe2x80x9cprecursorxe2x80x9d form, therefore, until cleavage of all or part of the antigen. A further masking may also be provided by molecules capable of binding to or otherwise associating with the B cell surface molecule binding domains of the B cell surface molecule binding portion. Examples of suitable masking molecules include but are not limited to the variable portion of a kappa or lambda light chain of an immunoglobulin.\nAlternatively, a fragment comprising a variable heavy and light chain (Fv) may be employed which is preferably but not exclusively a single chain (sc) and/or disulfide stabilized (ds). The nucleotide and corresponding amino acid sequences for the variable portion of kappa light chain are shown in SEQ ID NO:10 and SEQ ID NO:11, respectively.\nNotwithstanding that the abovementioned immunoglobulin portions are usefull molecules for blocking B cell surface immunoglobulin binding domains, other molecules may also be used. For example, natural product screening would very readily identify molecules from natural sources such as coral, soil, plants, ocean beds, marine invertebrates or from other convenient sources which would bind to immunoglobulin binding domains of L protein or other B cell surface binding molecules.\nIn a particularly preferred embodiment, the recombinant or synthetic growth factor precursor substantially prevents binding of X1 and X3 to their cognate B cell surface molecules thereby preventing B cell activation by having immunoglobulin peptides or chemical equivalents thereof or other B cell surface molecule blocking reagents linked, fused or otherwise associated with the growth factor to facilitate masking of the B cell activating effects of the growth factor. In a particularly preferred embodiment, the precursor comprises an antigen to which a catalytic antibody is sought and portions capable or masking the B cell surface molecule binding domains of the B cell surface molecule binding portion or the growth factor. The precursor may contain variable kappa or lambda light chain domains or an sc-ds-Fv molecule or masking molecules detected following natural product screening.\nGenerally, the immunoglobulin molecules which bind to the B cell surface immunoglobulin binding portion as a growth factor are linked to the N-terminal and C-terminal portions respectively of the antigen flanking the growth factor. For example, one particularly preferred embodiment of the present invention provides a growth factor precursor comprising the structure:\nIxe2x80x2AX1X2X3AIxe2x80x3\nwherein:\nIxe2x80x2 and Ixe2x80x3 are optionally present and may be the same or different and each is a blocking reagent for X1 and X3 such as a kappa or lambda light chain or a sc-ds-Fv molecule;\nA is the target antigen for which a catalytic antibody is sought;\nX1 and X3 are B cell surface molecule binding entities; and\nX2 is an entity providing T cell dependent help to a B cell.\nWherein a catalytic antibody on the surface of said B cell is capable of cleaving all or part of A from said recombinant or synthetic molecule resulting in the molecule [Axe2x80x2]X1X2X3[Axe2x80x2] wherein Axe2x80x2 is optionally present and is a portion of A after cleavage with the catalytic antibody wherein said resulting molecule is capable of inducing T cell dependent B cell mitogenesis of the B cell to which X1 and X3 bind.\nThe molecular components of Ixe2x80x2AX1X2X3AIxe2x80x3 may be in any sequence order.\nThe present invention extends to X1 and X3 being any B cell surface molecule binding molecules and components Ixe2x80x2 and IIxe2x80x3 block or mask the binding of such molecules to the B cell surface molecules.\nIn another embodiment the Ixe2x80x2A X1X2X3A Ixe2x80x3 molecule or part thereof may be in multimeric form, This is particularly the case when all or part of the molecule includes a multimerisation component (M) such as but not limited to the signal peptide of ompA. The monomeric units may be bound or otherwise associated together by any number of binding means such as contemplated above including covalent bonding, salt bridges, disulfide bridges and hydrophobic interactions amongst many others. Depending on the extent of multimerisation, this may impair the masking ability of B cell surface molecule binding domains of the growth factor and some antigen-independent clonal expansion may occur. This may not be too disadvantageous where there is at least some catalytic antibody dependent B cell mitogenesis.\nAccording to this embodiment, there is provided a growth factor precursor comprising the structure:\n[Ixe2x80x2AX1[X2xe2x80x2]o[X2X3AIxe2x80x3]n]m\nwherein:\nIxe2x80x2 and Ixe2x80x3 may be the same or different and each is a blocking reagent for X1 and X3 such as a kappa or lambda light chain or an sc-ds-Fv;\nA is the target antigen for which a catalytic antibody is sought;\nX1 and X3 are B cell surface molecule binding entities;\nX2 and X2xe2x80x2 may be the same or different and each is an entity capable or providing T cell dependent help for a B cell;\no may be 0or 1;\nn indicates the multimeric nature of the component in parentheses and may be 0, 1 or  greater than 1;\nm indicates the multimeric nature of the component in parenthesis and may be 1 or  greater than 1.\nPreferably, n and m are each from about 1 to about 10,000 more preferably from about 1 to about 1,000 and still more preferably from about 1 to about 200.\nPreferably, if n is 0, then o is 1.\nIn alternative embodiments, the growth factor precursor comprises the structure [[Ixe2x80x2AX2X3]n[X2xe2x80x2]o[X1AIxe2x80x3]m or [[Ixe2x80x2AX1[X2xe2x80x2]o]n[X2X3AIxe2x80x3]m].\nThe exact number ascribed to n and m may not be ascertainable but the multimeric nature identified functionally or physically by size (eg. determined using HPLC or PAGE).\nExamples of A include but are not limited to those hereinbefore described.\nThe present invention extends to the substitution of a non-antigen cleavage site such as a protease sensitive peptide which provides a useful model for designing growth factor molecules.\nAn exemplified growth actor precursor of the present invention is referred to herein as xe2x80x9cCATABxe2x80x9d and comprises a TNF flanked LHL with the variable region from a kappa and/or lambda light chain further masking the B surface immunoglobulin binding domain of the L molecules.\nIn another embodiment, the growth factor precursor mimics a CATAB being activatable by means other than or in addition to a catalytic antibody. For example, an enzyme sensitive molecule may be the xe2x80x9cantigenxe2x80x9d. In one embodiment, the CATAB mimic is referred to herein as xe2x80x9cTLHLxe2x80x9d and comprises a variable kappa light chain linked to the N-terminus of an amino acid linker sequence comprising the TEV protease cleavage site which is in turn linked to LHL. The nucleotide and amino acid sequences for TLHL are shown in SEQ ID NO:5 and SEQ ID NO:6. In another embodiment, the TLHL further comprises another amino acid linker sequence comprising the TEV protease cleavage site and a variable kappa light chain linked to the C-terminus of the latter amino acid linker sequence and is referred to herein as CATAB-TEV. The nucleotide sequence and corresponding amino acid sequence of CATAB is shown in SEQ ID NO:3 and SEQ ID NO:4, respectively.\nThe present invention extends to a growth factor precursor comprising a sequence of amino acids selected from:\n(i) the amino acid sequence set forth in SEQ ID NO:4 or having at least 60% similarity thereto (excluding the amino acid sequence of an antigen);\n(ii) an amino acid sequence encoded by the nucleotide sequence set forth in SEQ ID NO:3 or a nucleotide sequence having at least 60% similarity to the nucleotide sequence set forth in SEQ ID NO:3;\n(iii) an amino acid sequence encoded by a nucleotide sequence capable of hybridizing to SEQ ID NO:3 under low stringency conditions wherein catalytic properties of said growth factor precursor are capable of inducing B cell mitogenesis.\nThe variable kappa light chain portions may conveniently be derived from the Bence Jones protein LEN, which has been shown to bind to certain immunoglobulin binding proteins such as L. Alternatively, an Fv may be used such as a sc-ds-Fv.\nUpon cleavage of the growth factor precursor CATAB by a catalytic antibody recognising the antigen (for example, a TNF peptide portion), the covalent linkage between the L domains and the variable kappa light chains is broken. The blocking variable kappa light chains will dissociate from the L domains due to the relatively low affinity (xcx9c10xe2x88x927M) of individual domains for each other. This will release the mature growth factor LHL which can bind to and crosslink the surface immunoglobulin with an affinity of 10xe2x88x929M. A similar mechanism operates where molecules other than variable kappa light chains, TNF, L and H are employed.\nThis B cell surface immunoglobulin crosslinking activates the B cell and induces the down stream events that lead to catalytic maturation and secretion of the catalytic antibody. Catalytic antibodies can be detected in the serum using any number of procedures such as ELISA based assays and catalytic B cells may be recovered with standard hybridoma technology. Where the catalytic antibodies are from non-human animals, these can be humanised by recombinant DNA technology and produced for therapeutical applications in humans. Alternatively, the antibodies may be generated in a xe2x80x9chumanizedxe2x80x9d animal such as a humanized mouse which is transgenic for the human Ig loci.\nThe present invention contemplates derivatives of the growth factor and/or its precursor. A derivative includes a mutant, part, fragment, portion, homologue or analogue of the growth factor and/or precursor or any components thereof. Derivatives to amino acid sequences include single or multiple amino acid substitutions, deletions and/or additions. For example, derivatives of SEQ ID NO:2 or SEQ ED NO:4 include sequences having at least 60% similarity thereto (excluding amino acid sequence of antigen) or which are encoded by a nucleotide sequence capable of hybridizing in SEQ ID NO:1 or SEQ ID NO:3 under low stringency conditions.\nParticularly useful derivatives include chemical analogues of the growth factor and/or its precursor or components thereof. Such chemical analogues may be useful in stabilizing the molecule for therapeutic or diagnostic use.\nAnalogues of the growth factor precursor contemplated herein include, but are not limited to, modification to side chains, incorporating of unnatural amino acids and/or their derivatives during peptide, polypeptide or protein synthesis and the use of crosslinkers and other methods which impose conformational constraints on the proteinaceous molecule or their analogues.\nExamples of side chain modifications contemplated by the present invention include modifications of amino groups such as by reductive alkylation by reaction with an aldehyde followed by reduction with NaBH4; amidination with methylacetimidate; acylation with acetic anhydride; carbamoylation of amino groups with cyanate; trinitrobenzylation of amino groups with 2, 4, 6-trinitrobenzene sulphonic acid (TNBS); acylation of amino groups with succinic anhydride and tetrahydrophthalic anhydride; and pyridoxylation of lysine with pyridoxal-5-phosphate followed by reduction with NaBH4.\nThe guanidine group of arginine residues may be modified by the formation of heterocyclic condensation products with reagents such as 2,3-butanedione, phenylglyoxal and glyoxal.\nThe carboxyl group may be modified by carbodiimide activation via O-acylisourea formation followed by subsequent derivitisation, for example, to a corresponding amide.\nSulphydryl groups may be modified by methods such as carboxymethylation with iodoacetic acid or iodoacetamide; performic acid oxidation to cysteic acid; formation of a mixed disulphides with other thiol compounds; reaction with maleimide, maleic anhydride or other substituted maleimide; formation of mercurial derivatives using 4-chloromercuribenzoate, 4-chloromercuriphenylsulphonic acid, phenylmercury chloride, 2-chloromercuri-4-nitrophenol and other mercurials; carbamoylation with cyanate at alkaline pH.\nTryptophan residues may be modified by, for example, oxidation with N-bromosuccinimide or alkylation of the indole ring with 2-hydroxy-5-nitrobenzyl bromide or sulphenyl halides. Tyrosine residues on the other hand, may be altered by nitration with tetranitromethane to form a 3-nitrotyrosine derivative.\nModification of the imidazole ring of a histidine residue may be accomplished by alkylation with iodoacetic acid derivatives or N-carbethoxylation with diethylpyrocarbonate.\nExamples of incorporating unnatural amino acids and derivatives during peptide synthesis include, but are not limited to, use of norleucine, 4-amino butyric acid, 4-amino-3-hydroxy-5-phenylpentanoic acid, 6-aminohexanoic acid, t-butylglycine, norvaline, phenylglycine, ornithine, sarcosine, 4-amino-3-hydroxy-6-methylheptanoic acid, 2-thienyl alanine and/or D-isomers of amino acids. A list of unnatural amino acid, contemplated herein is shown in Table 1.\nCrosslinkers can be used, for example, to stabilise 3D conformations, using homo-bifunctional crosslinkers such as the bifunctional imido esters having (CH2)n spacer groups with n=1 to n=6, glutaraldehyde, N-hydroxysuccinimide esters and hetero-bifunctional reagents which usually contain an amino-reactive moiety such as N-hydroxysuccinimide and another group specific-reactive moiety such as maleimido or dithio moiety (SH) or carbodiimide (COOH). In addition, peptides can be conformationally constrained by, for example, incorporation of Cxcex1 and Nxcex1-methylamino acids, introduction of double bonds between Cxcex1 and Cxcex2 atoms of amino acids and the formation of cyclic peptides or analogues by introducing covalent bonds such as forming an amide bond between the N and C termini, between two side chains or between a side chain and the N or C terminus.\nThe present invention further contemplates chemical analogues of the growth factor and/or its precursor capable of acting as antagonists or agonists of same. These may be useful in controlling the immunological response. Chemical analogues may not necessarily be derived from the growth factor precursor but may share certain conformational similarities. Alternatively, chemical analogues may be specifically designed to mimic certain physiochemical properties of the growth factor precursor. Chemical analogues may be chemically synthesised or may be detected following, for example, natural product screening of, for example, coral, soil, plants, microorganisms, marine invertebrates or seabeds.\nOther derivatives contemplated by the present invention include a range of glycosylation variants from a completely unglycosylated molecule to a modified glycosylated molecule. Altered glycosylation patterns may result from expression of recombinant molecules in different host cells.\nThe growth factor and growth factor precursor of the present invention may be produced by recombinant DNA means or by chemical synthetic processes. With respect to the former this aspect of the present invention provides a nucleic acid molecule comprising a sequence of nucleotides encoding a B cell immunoglobulin binding peptide or polypeptide and a T cell epitope.\nGenerally, the nucleic acid molecule encodes a fusion molecule comprising a B cell immunoglobulin binding peptide or polypeptide and a T cell epitope. In this regard, the nucleic acid molecule would generally comprise a promoter region with the B cell immunoglobulin binding peptide or polypeptide encoding portion operably linked thereto and the T cell epitope portion operably linked to the 3xe2x80x2 end thereof. Expression of the nucleic acid molecule of the present invention leads to synthesis of a fusion molecule. In a preferred embodiment the nucleic acid molecule further encodes an antigenic portion and optionally an immunoglobulin portion such as a variable kappa light chain. In some circumstances, purification of one or more recombinant components results in all or a portion of the signal peptide for the ompA or other carrier molecule for the microbial host in which the recombinant molecule is made. This may facilitate multimerisation. In a particularly preferred embodiment of the present invention, the nucleic acid molecule comprises the nucleotide sequence set forth in SEQ ID NO:1 or SEQ ID NO:3 or is a sequence capable of hybridizing thereto under low stringency conditions wherein said nucleotide sequence encodes a growth factor precursor as hereinbefore described.\nStill a further aspect of the present invention extends to a method for producing catalytic antibodies to a specific antigen, said method comprising administering to an animal an effective amount of a growth factor precursor comprising an antigen capable of interacting with a B cell bound catalytic antibody said antigen linked to or otherwise associate with a B cell surface molecule binding portion and a portion capable of providing T cell dependent help to a B cell. The growth factor precursor may also comprise a B cell surface molecule binding portion masking entity such as a variable kappa light chain linked to the antigen.\nAlternatively, the growth factor precursor may comprise a B cell surface molecule binding portion in multimeric form linked to an antigen for which a target antibody is sought. The portion providing T cell dependent help is preferably a T cell epitope and is preferably part of the precursor. However, it may be a separate entity administered simultaneously or sequentially to an animal.\nThe present invention also provides catalytic antibodies produced by the above method. The catalytic antibodies are in effect specific to or directed against A of the above formulae. Such catalytic antibodies, therefore, may be directed to any antigen such as but not limited to a cytokine such as but not limited to tumor necrosis factor (TNF), an interleukin (IL) such as IL-1 to IL-15, interferons (IFN) such as IFNxcex1, IFNxcex2 or IFNxcex3, colony-stimulating factors (CSF) such as granulocyte colony-stimulating factor (G-CSF), ganulocyte-macrophase colony-stimulation factor (GM-CSF), blood factors such as Factor VIII, erythropoietin and haemopoietin, cancer antigens, docking receptors from pathogenic viruses such as HIV, influenza virus or a hepatitis virus (eg. HEP A, HEP B, HEP C or HEP E) and amyloid plaques such as in Alzheimer\"\"s disease patients or myeloma patients.\nThe catalytic antibodies of the present invention have particular therapeutic and diagnostic uses especially in relation to mammalian and more particularly human disease conditions.\nAccordingly, the present invention contemplates a pharmaceutical composition comprising a growth factor precursor or a derivative thereof and optionally a modulator of growth factor precursor activity and one or more pharmaceutically acceptable carriers and/or diluents. More particularly, the pharmaceutical composition comprises catalytic antibodies generated by the growth factor precursor of the present invention. These components are hereinafter referred to as the xe2x80x9cactive ingredientsxe2x80x9d.\nThe pharmaceutical forms suitable for injectable use include sterile aqueous solutions (where water soluble) or dispersions and sterile powders for the extemporaneous preparation of sterile injectable solutions or dispersion or may be in the form of a cream or other form suitable for topical application. It must be stable under the conditions of manufacture and storage and must be preserved against the contaminating action of microorganisms such as bacteria and fungi. The carrier can be a solvent or dispersion medium containing, for example, water, ethanol, polyol (for example, glycerol, propylene glycol and liquid polyethylene glycol, and the like), suitable mixtures thereof, and vegetable oils. The proper fluidity can be maintained, for example, by the use of a coating such as licithin, by the maintenance of the required particle size in the case of dispersion and by the use of superfactants. The preventions of the action of microorganisms can be brought about by various antibacterial and antingal agents, for example, parabens, chlorobutanol, phenol, sorbic acid, thirmerosal and the like. In many cases, it will be preferable to include isotonic agents, for example, sugars or sodium chloride. Prolonged absorption of the injectable compositions can be brought about by the use in the compositions of agents delaying absorption, for example, aluminum monostearate and gelatin.\nSterile injectable solutions are prepared by incorporating the active compounds in the required amount in the appropriate solvent with various of the other ingredients enumerated above, as required, followed by filtered sterilization. Generally, dispersions are prepared by incorporating the various sterilized active ingredient into a sterile vehicle which contains the basic dispersion medium and the required other ingredients from those enumerated above. In the case of sterile powders for the preparation of sterile injectable solutions, the preferred methods of preparation are vacuum drying and the freeze-drying technique which yield a powder of the active ingredient plus any additional desired ingredient from previously sterile-filtered solution thereof.\nWhen the active ingredients are suitably protected they may be orally administered, for example, with an inert diluent or with an assimilable edible carrier, or it may be enclosed in hard or soft shell gelatin capsule, or it may be compressed into tablets, or it may be incorporated directly with the food of the diet. For oral therapeutic administration, the active compound may be incorporated with excipients and used in the form of ingestible tablets, buccal tablets, troches, capsules, elixirs, suspensions, syrups, wafers, and the like. Such compositions and preparations should contain at least 1% by weight of active compound. The percentage of the compositions and preparations may, of course, be varied and may conveniently be between about 5 to about 80% of the weight of the unit. The amount of active compound in such therapeutically useful compositions in such that a suitable dosage will be obtained. Preferred compositions or preparations according to the present invention are prepared so that an oral dosage unit form contains between about 0.1 ng and 2000 mg of active compound, preferably between about 0.1 xcexcg and 1500 mg and more preferably between about 1 xcexcg and 100 mg.\nThe tablets, troches, pills, capsules and the like may also contain the components as listed hereafter: a binder such as gum, acacia, corn starch or gelatin; excipients such as dicalcium phosphate; a disintegrating agent such as corn starch, potato starch, alginic acid and the like; a lubricant such as magnesium stearate; and a sweetening agent such a sucrose, lactose or saccharin may be added or a flavouring agent such as peppermint, oil of wintergreen, or cherry flavouring. When the dosage unit form is a capsule, it may contain, in addition to materials of the above type, a liquid carrier. Various other materials may be present as coatings or to otherwise modify the physical form of the dosage unit. For instance, tablets, pills, or capsules may be coated with shellac, sugar or both. A syrup or elixir may contain the active compound, sucrose as a sweetening agent, methyl and propylparabens as preservatives, a dye and flavouring such as cherry, orange or mango. Of course, any material used in preparing any dosage unit form should be pharmaceutically purer substantially non-toxic in the amounts employed. In addition, the active compound(s) may be incorporated into sustained-release preparations and formulations.\nThe present invention also extends to forms suitable for topical application such as creams, lotions and gels.\nPharmaceutically acceptable carriers and/or diluents include any and all solvents, dispersion media, coatings, antibacterial and antifungal agents, isotonic and absorption delaying agents and the like. The use of such media and agents for pharmaceutical active substances is well known in the art. Except insofar as any conventional media or agent is incompatible with the active ingredient, use thereof in the therapeutic compositions is contemplated. Supplementary active ingredients can also be incorporated into the compositions. These may include immune potentiating molecules, multimer facilitating molecules and pharmaceutically active molecules chosen on the disease conditions being treated.\nIt is especially advantageous to formulate parenteral compositions in dosage unit form for ease of administration and uniformity of dosage. Dosage unit form as used herein refers to physically discrete units suited as unitary dosages for the mammalian subjects to be treated; each unit containing a predetermined quantity of active material calculated to produce the desired therapeutic effect in association with the required pharmaceutical carrier. The specification for the novel dosage unit forms of the invention are dictated by and directly dependent on (a) the unique characteristics of the active material and the particular therapeutic effect to be achieved, and (b) the limitations inherent in the art of compounding such an active material for the treatment of disease in living subjects having a diseased condition in which bodily health is impaired as herein disclosed in detail.\nThe principal active ingredient is compounded for convenient and effective administration in effective amounts with a suitable pharmaceutically acceptable carrier in dosage unit form as hereinbefore disclosed. A unit dosage form can, for example, contain the principal active compound in amounts ranging from 0.1 ng to about 200 mg, more preferably ranging from 0.1 xcexcg and 1500 mg and even more preferably ranging between 1 xcexcg and 1000 mg. Expressed in proportions the active compound is generally present in from about 0.5 xcexcg to about 2000 mg/ml of carrier. In the case of compositions containing supplementary active ingredients, the dosages are determined by reference to the usual dose and manner of administration of the said ingredients.\nStill another aspect of the present invention is directed to antibodies to the growth factor precursor and its derivatives. Such antibodies may be monoclonal or polyclonal and are independent to the catalytic antibodies selected by the precursor. The (non-catalytic) antibodies to recombinant or synthetic the growth factor precursor or its derivatives of the present invention may be useful as therapeutic agents but are particularly useful as diagnostic agents. Antibodies may also be generated to the catalytic antibodies generated by the growth factor precursors. All these antibodies have particular application in diagnostic assays for the growth factor or catalytic antibody inducer thereof.\nFor example, specific antibodies can be used to screen for catalytic antibodies. The latter would be important, for example, as a means for screening for levels of these antibodies in a biological fluid or for purifying the catalytic antibodies. Techniques for the assays contemplated herein are known in the art and include, for example, sandwich assays and ELISA.\nIt is within the scope of this invention to include any second antibodies (monoclonal, polyclonal or fragments of antibodies or synthetic antibodies) directed to the antibodies discussed above. Both the first and second antibodies may be used in detection assays or a fist antibody may be used with a commercially available anti-immunoglobulin antibody.\nBoth polyclonal and monoclonal antibodies are obtainable by immunization with the enzyme or protein and either type is utilizable for immunoassays. The methods of obtaining both types of sera are well known in the art. Polyclonal sera are less preferred but are relatively easily prepared by injection of a suitable laboratory animal with an effective amount of antigen, or antigenic parts thereof, collecting serum from the animal, and isolating specific sera by any of the known immunoadsorbent techniques. Although antibodies produced by this method are utilizable in virtually any type of immunoassay, they are generally less favoured because of the potential heterogeneity of the product.\nThe use of monoclonal antibodies in an immunoassay is particularly preferred because of the ability to produce them in large quantities and the homogeneity of the product. The preparation of hybridoma cell lines for monoclonal antibody production derived by fusing an immortal cell line and lymphocytes sensitized against the immunogenic preparation can be done by techniques which are well known to those who are skilled in the art.\nAnother aspect of the present invention contemplates a method for detecting an antigen in a biological sample from a subject said method comprising contacting said biological sample with an antibody specific, for said antigen or its derivatives or homologues for a time and under conditions sufficient for an antibody-antigen complex to form, and then detecting said complex. In this context the xe2x80x9cantigenxe2x80x9d may be a growth factor, its precursor, a component thereof or a catalytic antibody induced thereby.\nThe presence of antigen may be accomplished in a number of ways such as by Western blotting and ELISA procedures. A wide range of immunoassay techniques are available as can be seen by reference to U.S. Pat. Nos. 4,016,043, 4,424,279 and 4,018,653. These, of course, includes both single-site and two-site or xe2x80x9csandwichxe2x80x9d assays of the non-competitive types, as well as in the traditional competitive binding assays. These assays also include direct binding of a labelled antibody to a target.\nSandwich assays are among the most useful and commonly used assays and are favoured for use in the present invention. A number of variations of the sandwich assay technique exist, and all are intended to be encompassed by the present invention. Briefly, in a typical forward assay, an unlabelled antibody is immobilized on a solid substrate and the sample to be tested brought into contact with the bound molecule. After a suitable period of incubation, for a period of time sufficient to allow formation of an antibody-antigen complex, a second antibody specific to the antigen, labelled with a reporter molecule capable of producing a detectable signal is then added and incubated, allowing time sufficient for the formation of another complex of antibody-antigen-labelled antibody. Any unreacted material is washed away, and the presence of the antigen is determined by observation of a signal produced by the reporter molecule. The results may either be qualitative, by simple observation of the visible signal, or may be quantitated by comparing with a control sample containing known amounts of hapten. Variations on the forward assay include a simultaneous assay, in which both sample and labelled antibody are added simultaneously to the bound antibody. These techniques are well known to those skilled in the art, including any minor variations as will be readily apparent In accordance with the present invention the sample is one which might contain an antigen including cell extract, supernatant fluid, tissue biopsy or possibly serum, saliva, mucosal secretions, lymph, tissue fluid and respiratory fluid. The sample is, therefore, generally a biological sample comprising biological fluid but also extends to fermentation fluid and supernatant fluid such as from a cell culture.\nIn the typical forward sandwich assay, a first antibody having specificity for the antigen or antigenic parts thereof, is either covalently or passively bound to a solid surface. The solid surface is typically glass or a polymer, the most commonly used polymers being cellulose, polyacrylamide, nylon, polystyrene, polyvinyl chloride or polypropylene. The solid supports may be in the form of tubes, beads, discs of microplates, or any other surface suitable for conducting an immunoassay. The binding processes are well-known in the art and generally consist of cross-lining covalently binding or physically adsorbing, the polymer-antibody complex is washed in preparation for the test sample. An aliquot of the sample to be tested is then added to the solid phase complex and incubated for a period of time sufficient (e.g. 2-40 minutes) and under suitable conditions (e.g. 25xc2x0 C.) to allow binding of any subunit present in the antibody. Following the incubation period, the antibody subunit solid phase is washed and dried and incubated with a second antibody specific for a portion of the hapten. The second antibody is linked to a reporter molecule which is used to indicate the binding of the second antibody to the hapten.\nAn alternative method involves immobilizing the target molecules in the biological sample and then exposing the immobilized target to specific antibody which may or may not be labelled with a reporter molecule. Depending on the amount of target and the strength of the reporter molecule signal, a bound target may be detectable by direct labelling with the antibody. Alternatively, a second labelled antibody, specific to the first antibody is exposed to the target-first antibody complex to form a target-first antibody-second antibody tertiary complex. The complex is detected by the signal emitted by the reporter molecule.\nBy xe2x80x9creporter moleculexe2x80x9d as used in the present specification, is meant a molecule which, by its chemical nature, provides an analytically identifiable signal which allows the detection of antigen-bound antibody. Detection may be either qualitative or quantitative. The most commonly used reporter molecules in this type of assay are either enzymes, fluorophores or radionuclide containing molecules (i.e. radioisotopes) and chemiluminescent molecules. In the case of an enzyme immunoassay, an enzyme is conjugated to the second antibody, generally by means of glutaraldehyde or periodate. As will be readily recognized, however, a wide variety of different conjugation techniques exist, which are readily available to the skilled artisan. Commonly used enzymes include horseradish peroxidase, glucose oxidase, beta-galactosidase and alkaline phosphatase, amongst others. The substrates to be used with the specific enzymes are generally chosen for the production, upon hydrolysis by the corresponding enzyme, of a detectable colour change. Examples of suitable enzymes include alkaline phosphatase and peroxidase. It is also possible to employ fluorogenic substrates, which yield a fluorescent product rather than the chromogenic substrates noted above. In all cases, the enzyme-labelled antibody is added to the first antibody hapten complex, allowed to bind, and then the excess reagent is washed away. A solution containing the appropriate substrate is then added to the complex of antibody-antigen-antibody. The substrate will react with the enzyme linked to the second antibody, giving a qualitative visual signal, which may be further quantitated, usually spectrphotometrically, to give an indication of the amount of hapten which was present in the sample. xe2x80x9cReporter moleculexe2x80x9d also extends to use of cell agglutination or inhibition of agglutination such as red blood cells on latex beads, and the like.\nAlternately, fluorescent compounds, such as fluorescein and rhodamine, may be chemically coupled to antibodies without altering their binding capacity. When activated by illumination with light of a particular wavelength, the fluorochrome-labelled antibody adsorbs the light energy, inducing a state to excitability in the molecule, followed by emission of the light at a characteristic colour visually detectable with a light microscope. As in the EIA, the fluorescent labelled antibody is allowed to bind to the first antibody-hapten complex. After washing off the unbound reagent, the remaining tertiary complex is then exposed to the light of the appropriate wavelength the fluorescence observed indicates the presence of the hapten of interest. Immunofluorescene and EIA techniques are both very well established in the art and are particularly preferred for the present method. However, other reporter molecules, such as radioisotope, chemiluminescent or bioluminescent molecules, may also be employed.\nThe present invention may use any number of means to clone genetic sequences encoding catalytic antibodies. For example, a phage display library potentially capable of expressing a catalytic antibody on the phage surface may be used to screen for catalysis of defined antigens.\nThe present invention further contemplates the use of the products of catalysis of a growth factor precursor to induce B cell mitogenesis to generate catalytic antibodies to a specific antigen.\nMore particularly, the present invention contemplates the use of a growth factor precursor comprising an antigen to which a catalytic antibody is sought linked, fused or otherwise associated to a B cell surface molecule binding portion in the induction of B cell mitogenesis following catalytic cleavage of all or part of said antigen.\nStill another embodiment of the present invention contemplates the use of an antigen linked, fused or otherwise associate to a B cell surface molecule binding portion in the manufacture of a growth factor precursor to induce B cell mitosis following catalytic cleavage of all or part of said antigen."}
{"text": "1. Field of the Invention\nThe invention relates to a power amplifier, and in particular to a power amplifier circuit and method thereof.\n2. Description of the Related Art\nSpectrum usage in technology and telecommunication industries is typically governed by a spectrum regulator such as the Federal Communications Commission (FCC) for transmission power over a particular spectrum to manage interference and spectrum shortage.\nIn practice, a sharp change in signal envelope amplitude usually exerts frequency components outside the regulated power spectrum. This phenomenon is known as “spectral spatter”, referring to the frequency “spattering” beyond its intended spectrum boundaries. Spectral spatter is usually not desirable, since it not only produces undesirable interference on adjacent spectra, but also is energy inefficient.\nVarious efforts have been made to reduce the spectral spatter, such as applying a filter to remove excess frequency components. However, a filter may not effectively remove undesired frequency components, and attenuates the magnitude of spectrum in use. Another attempt focuses on power amplifier circuit with ramp control. This approach retards the amplitude variation of signal envelope. Whenever the ramp circuit experiences an abrupt change at its input, it delivers a gradually inclining signal transition to the power amplifier. Consequently the power amplifier follows the inclining signal transition to output a progressive slanted signal envelope, whereby the excess frequency components are suppressed.\nFIG. 1 is a block diagram of a conventional power amplifier circuit with ramp control, comprising operational amplifier (OA) 10, MOS 12, and power amplifier (PA) 14. The output of OA 10 is coupled to the gate of MOS 12, the inverting input of OA 10 is coupled to the drain of MOS 12 and a supply power input of PA 14.\nWhen the conventional power amplifier circuit is required to amplify input signal Sin1 to output signal Sout1, the non-inverting input of OA 10 receives enable signal Sen1 from logic “low” to logic “high”, such that the output of OA 10 generates a bias voltage to turn MOS 12 on, rendering the supply voltage at the supply voltage input of PA 14. The supply voltage is also fedback to the inverting input of OA 10, so that as the voltage difference between the inverting and non-inverting inputs of OA 10 decreases, the bias voltage decreases accordingly, leading to a ramp voltage to the supply voltage input of PA 14. Consequently PA 14 amplifies input signal Sin1 by a variable gain corresponding to the ramp voltage from the supply voltage input, and generates output voltage Sout1 with a ramp envelop.\nAlthough the conventional power amplifier circuit provides ramp control capability, it also presents inherent circuit deficiency. Firstly, PA 10 in the conventional power amplifier circuit increases circuit complexity. Secondly, power consumption in MOS 12 causes ineffective power utilization. Thirdly, the supply voltage to PA 14 cannot complete a full VCC swing since there is a finite voltage drop across MOS 12. Finally, the device size of MOS 12 has to be large to withstand a large driving current, consuming more circuit area and increasing manufacturing cost.\nIn addition, the power amplifier circuit is also required to control the initialization of the ramp bias signal. If the input signal arrives at the power amplifier after the ramp bias signal reaches a maximal or minimal value the spectral spatter can still be present in the output signal. As a result a programmable delay can be incorporated into the power amplifier circuit such that the ramp bias signal is synchronous with the input signal.\nThus it is desirable to have a power amplifier circuit capable of ramp control and programmable delay control."}
{"text": "Vanity mirrors which mount to sun visors have long been a popular accessory in passenger automobiles. Typically, the mirror is mounted on the upper surface of the sun visor, so that when the sun visor is swung downwardly to a sun shading position, the mirror is aligned with the face of the driver or passenger for uses such as personal grooming and the like. Certain of such vanity mirrors have been provided in packages with electric lighting which illuminates the users face for night time operation or to approximate ambient light conditions to be encountered by the user at his or her destination. One known type of lighted vanity mirror includes a pair of electric lamps disposed at opposite ends of the mirror and a pivotable cover for covering the mirror and lights when not in use. Circuitry for connection of the lights to the automobile's electrical system, including an on/off switch and often times a dimming control, is also provided in the module.\nIt has been found that such lighted vanity mirror modules can be improved upon in a number of respects. One potential area of improvement lies in the packaging of the electrical circuitry employed in the module. Traditionally, most illuminated vanity mirror modules comprise a shallow housing covered by an appearance bezel which usually supports the mirror, lid, and lenses for the electric lamps and the mirror. The electrical circuitry consisting of an on/off switch in some cases a dimming control, a pair of bulbs, and electrical conductors comprising metallic stampings or wires including the bulbs, switch and dimmer, are disposed in the housing beneath the bezel. Where the circuit conductors comprise metallic stampings, such stampings are often heat staked to, or molded within a base plate of the housing.\nSuch constructions are disadvantageous in a number of respects. For example, the housing and bezel tend to be heavy and bulky, thereby adding significantly to the weight and bulk of the visor on which they are attached. This renders the visor cumbersome to operate and dimensionally, quite thick thereby reducing passenger head room in the automobile. Molding the metallic stampings into the housing base plate when the base plate, adds significantly to the cost of molding the housing and risks unacceptable warpage of the base plate during manufacture.\nIn an effort to reduce the bulk and weight of the lighted mirror module, it is sometimes the practice to mount the electrical circuit components directly to a molded sun visor blade. While this construction eliminates the need for a housing and base plate, it limits the type of visors to which the module is adaptable and, therefore, may not be useful with the padded and upholstered visors found in current luxury automobiles. Moreover, if the lighted mirror module requires service, the entire visor must be serviced since that the module is an integral part of the visor and is not conveniently removable therefrom for servicing.\nAnother area of improvement in current lighted vanity mirror modules lies in the construction of the dimming control. Certain mirror modules employ resistors connected along the length thereof to a plurality of contacts formed into a metallic stamping and a movable wiper contact which rides along the contacts to adjust the amount of resistance in the lamp circuit and thus control the lamp intensity. Such a construction is found in U.S. patent application Ser. No. 07/357,652. While such a construction has proven to be effective and reliable, the metallic stamping employed into which the resistor contacts are employed, is rather complex in shape, somewhat costly, and therefore not as economical as may be desired. Other lighted mirror modules such as that disclosed in U.S. Pat. No. 4,879,637 employ specialized resistors comprising resistive coatings on a circuit board. However, such resistors can be prohibitively expensive and may require multiple contacts similar to those discussed hereinabove and, therefore, are also characterized by some significant cost and complexity.\nAccordingly, a lighted vanity mirror module characterized by enhanced compactness, simplicity of construction and economy of manufacture is highly desirable."}
{"text": "Portable terminals have become necessities of modern life for people of all ages. Thus, service providers and terminal manufacturers are competitively developing differentiated products and services.\nFor example, the portable terminal has developed into a multimedia device capable of providing various services such as phonebooks, games, Short Message Service (SMS), Electronic mail (E)-mail, morning calls, MPEG-1 Audio Layer 3 (MP3) player, scheduling, digital cameras, Multimedia Message Service (MMS), and wireless Internet.\nAdditionally, the aforementioned portable terminal provides a location information service which measures a location of the portable terminal, i.e., a location of a user, to inform the location. By using the location information service, the user can determine a user's current location, a bus stop around the current location, a path to a destination from the current location, and so forth. The location information service is one service frequently used by the user.\nThe location information service uses a Global Positioning System (GPS) in general. The GPS can obtain correct information in an outdoor zone where a satellite signal can be readily received, but cannot be used in an indoor zone where the satellite signal cannot be received.\nIn order to use the location information service in the indoor zone where the satellite signal cannot be received, a Wi-Fi Positioning System (WPS) technique that utilizes a Wi-Fi signal can be used. The WPS technique can be classified into a triangulation technique and a fingerprint technique. However, since Wireless Local Area Network (WLAN) signal strength is measured and used in both of the two techniques, at least 3 or more Access Points (APs) may be required to provide the location information service with high accuracy. Further, there is a problem in that AP information (e.g., Service Set IDentifier (SSID), Media Access Control (MAC) address, location, and the like) has to be managed in advance in a form of a database.\nInstead of the aforementioned method, a sensor may be installed in a lattice form in an indoor zone (floor, ceiling, and the like), and thereafter a user location may be recognized by using a difference in a signal transfer rate between an ultrasonic signal and a signal generated from the sensor, or the user location can be recognized by utilizing a plurality of cameras, a pressure sensor, an infrared ray, and a Radio Frequency IDentification (RFID). However, this method may result in great expenses to install the sensor.\nIn addition, when using a method in which a magnetic field sensor is utilized to recognize the user location, a magnetic field generator is used to implement the method. Therefore, this method also has a problem in that the installation cost increases similarly to the case of using the sensor."}
{"text": "Gene transfer to the alveolar epithelium is an attractive therapeutic approach for a number of acute and chronic acquired lung diseases, including pulmonary inflammation, pulmonary edema, acute lung injury (ALI), the acute respiratory distress syndrome (ARDS), and pulmonary fibrosis (1, 5, 15, 43). In addition, due to its large surface area and proximity to the vascular endothelium, the alveolar epithelium represents a desirable target for delivery of therapeutic genes encoding secreted proteins. However, in contrast to the large number of studies that have utilized a variety of vectors to achieve gene transfer to tracheal and bronchial epithelium in the upper airways, particularly in the context of gene therapy for cystic fibrosis (54), relatively few studies have examined gene transfer into the alveolar epithelium of the distal respiratory tract.\nNon-viral strategies for delivery of exogenous DNA to the lung have been limited by low efficiency of transduction (56). Of the various vectors that have been evaluated for gene therapy thus far, each exhibits characteristic disadvantages, and none has proven effective in achieving efficient, long-term expression in the distal respiratory tract. Due to the relative quiescence of the cells that constitute the alveolar epithelium, viral vectors must be able to efficiently transduce cells that are not actively dividing. In this regard, adenoviral vectors have been shown to effectively transduce the alveolar epithelium (15). However, use of these vectors in vivo has been limited by immune responses, which can be especially problematic in the lung alveoli due to the potential for inducing serious pulmonary inflammation. In any event, use of this non-integrating vector system would also require repeated viral administration to achieve long-term gene expression (12). Adeno-associated virus (AAV) vectors show episomal expression and eventual integration following cell division and have been used for gene delivery to the distal respiratory tract (16, 17); however, their use has been limited by low packaging capacity and difficulty obtaining high titer preparations (20). Repeated administration of AAV vectors has been limited by the development of neutralizing antibodies (8). Thus, investigation of alternative vectors, such as retroviruses, for gene delivery to the alveolar epithelium is warranted.\nTwo types of retrovirus vectors, murine leukemia virus-(MLV) and lentivirus-based, have been tested for their efficiency in transducing alveolar epithelial cells (“AEC”). Under normal in vivo conditions, the cells that constitute the alveolar epithelium undergo very low rates of proliferation (2, 4, 46, 47). The efficiency of standard MLV-based retroviruses for gene transfer to the adult alveolar epithelium is therefore predictably inefficient (14, 50). This limitation has been partially overcome by inducing cell proliferation with growth factors, but overall transduction efficiency is still quite low. Newer lentivirus vectors, however, have recently been shown to transduce several non-dividing cell types including neuronal cells, myocytes, liver cells and tracheal epithelial cells, with long term persistence of transgene expression (21, 23, 34, 40).\nPseudotyping of the lentivirus with different envelope proteins has expanded the range of host cells that can be transduced by lentivirus vectors, and has also allowed the virus to be easily concentrated to high titers, especially when pseudotyped with the vesicular stomatitis virus envelope glycoprotein (VSV-G) (38). Lentivirus vectors would therefore appear to be ideally suited for gene transfer to the relatively quiescent cells of the alveolar epithelium. However, the lack of a well-characterized virus receptor and the polarized nature of its distribution makes transfection efficiency in different cell types unpredictable.\nStudies to date have focused on lentivirus-based gene transfer to the tracheal or bronchial epithelium of the proximal airways. In polarized, well-differentiated airway epithelia, only minimal transduction by VSV-G pseudotyped vectors introduced from the apical surface (the only directly accessible surface in vivo) has been observed. The use of lentivirus vectors for transduction of adult alveolar epithelial cells in the distal respiratory tract has not been evaluated to date. In particular, whether the polarized cells that constitute the alveolar epithelium present a similar barrier to apical transduction by lentivirus has not been explored.\nBased on the 1984 paper by Fuller et al., it was widely believed in the art that receptors for VSV-G were exclusively located on the basolateral cell surface (41). Because only the apical surface is accessible in AEC in vivo, it was believed that it would be futile to use VSV-G pseudotyped lentivirus to transduce these cells.\nThe present invention overcomes these hurdles and provides a composition and method for highly efficient transduction of AEC from the apical surface using the lentivirus vector."}
{"text": "1. Field of the Invention\nThis invention relates broadly to the field of packaging and more specifically to a tamper evident container closure for packaging products such as beverages.\n2. Description of the Related Technology\nPlastic closures for containers such as beverage bottles are in widespread use throughout the world. In recent years, circumstances in many countries has resulted in an increase in demand for tamper evident features on such closures.\nAs is described in U.S. Pat. No. 5,392,938, the use of tamper evident security bands in plastic screw closures for bottles having a screw cap has been common for some time, but the use of such security bands on snap hinge closures has been much less common. In all of the snap hinge closures of which the inventors are aware, the body, the lid and the hinge of the closure is molded together in one piece and mechanical interlocking structure is formed in the body and/or the lid for causing the lid to become mechanically affixed to the body when the lid is pressed onto the body after manufacturing. This mechanical connection is designed to be defeated in a tamper evident manner, such as by rupturing of frangible bars or webs within the connection when the closure is first opened by a consumer.\nUnfortunately, in too many circumstances the mechanical connections described above are able to be defeated without rupturing the tamper evident bars or webs by a person who is determined to do so. This, of course, is unacceptable for many different reasons. A need exists for an improved snap hinge tamper evident closure and a method for making such a closure that is impossible to open without defeating the tamper evident structure of the closure."}
{"text": "1. Field of the Invention\nThe present invention relates to a load drive circuit and a display device, and in particular to a drive circuit and a display device successfully reduced in unnecessary radiation in operation of any display panel, which acts as a load, such as plasma display electro-luminescence display and liquid crystal display (LCD).\n2. Description of the Related Art\nFIG. 10 is a schematic drawing of a three electrode surface discharge AC-driven plasma display panel, and FIG. 11 is a sectional view for explaining the electrode structure of the plasma display panel shown in FIG. 10. In FIGS. 10 and 11, reference numeral 207 denotes a discharge cell (display cell), 210 denotes a rear glass substrate, 211 and 221 denote dielectric layers, 212 denotes a fluorescent material, 213 denotes a diaphragm, 214 denotes address electrodes (A1 to Ad), 220 denotes a front glass substrate, and 222 denotes X electrodes (X1 to XL) or Y electrodes (Y1 to YL). Reference symbol Ca denotes capacitance between the adjacent address electrodes, and Cg denotes capacitance between the opposing address electrodes (X electrode and Y electrode).\nA plasma display panel 201 comprises two glass substrates of the rear glass substrate 210 and the front glass substrate 220, and the front glass substrate 220 has the X electrodes (X1, X2, . . . , XL) and the Y electrodes (scanning electrodes: Y1, Y2, to YL) disposed thereon as sustaining electrodes (including bus electrode and transparent electrode).\nThe rear glass substrate 210 has the address electrodes (A1, A2, . . . , Ad) 214 disposed thereon so as to cross normal to the sustaining electrodes (X electrodes and Y electrodes) 222, and the display cell 207 causing discharge light emission with the aid of these electrodes is formed in each area which falls between the X electrode and Y electrode having the same number (Y1-X1, Y2-X2, . . . ), and on the intersection with the address electrode.\nFIG. 12 is a block diagram showing an entire configuration of a plasma display device using the plasma display panel shown in FIG. 10, and more specifically showing an essential portion of a drive circuit for the display panel.\nAs shown in FIG. 12, the three electrode surface discharge, AC-driven plasma display device comprises the display panel 201, a control circuit 205 for generating a control signal for controlling the drive circuit of the display panel in response to an externally-input interface signal, an X common driver (X electrode drive circuit) 206 for driving the panel electrodes using the control signal from the control circuit 205, a scanning electrode drive circuit (scan driver) 203 and Y common driver 204, and an address electrode drive circuit (address driver) 202.\nThe X common driver 206 generates sustaining voltage pulse, the Y common driver 204 similarly generates sustaining voltage pulse, and the scan driver 203 effects scanning by independently driving the individual scan electrodes (Y1 to YL). The address driver 202 applies an address voltage pulse to the individual address electrodes (A1 to Ad) corresponding to display data.\nThe control circuit 205 has a display data control section 251 for supplying an address control signal to the address driver 202 upon receiving of clock CLK and display data DATA, a scan driver control section 253 for controlling the scan driver 203 upon receiving of vertical synchronizing signal Vsync and horizontal synchronizing signal Hsync, and a common driver control section 254 for controlling common drivers (X common driver 206 and Y common driver 204). The display data control section 251 has a frame memory 252.\nFIG. 13 is a drawing showing an exemplary drive waveform of the plasma display device shown in FIG. 12, and shows outlines of voltage waveforms applied to the individual electrodes in the all-write period (AW), all-erasure period (AE), address period (ADD) and sustaining period (sustained discharge period: SUS).\nIn FIG. 13, drive periods directly related to image display are the address period ADD and sustaining period SUS, where a pixel to be displayed is selected in the address period ADD, and the selected pixel is then kept in an emission state so as to allow image display in a predetermined brightness. It is to be noted that FIG. 13 shows drive waveforms of the individual sub-frames for the case where one frame is composed of a plurality of subframes (sub-field).\nFirst, in the address period ADD, all of the Y electrodes (Y1 to YL), which are scanning electrodes, are applied with an intermediate potential −Vmy en bloc, and then sequentially applied with a scanning voltage pulse of −Vy level as being switched from −Vmy. In this process, pixels on the individual scanning lines can be selected by applying an address voltage pulse of +Va level to the individual address electrodes (electrode A: A1 to Ad) in synchronization with the application of the scanning pulses to the individual Y electrodes.\nIn the succeeding sustaining period SUS, all of the scanning electrodes (Y1 to YL) and X electrodes (X1 to XL) are alternately applied with a sustaining voltage pulse of +Vs level, so as to induce sustained emission at the previously selected pixels, and to allow display at a predetermined luminance through such successive application. It is also made possible to effect gradation display by controlling the number of times of light emission by combining basic operations expressed by such series of operation waveforms.\nThe all-write period AW is provided for applying a write voltage pulse to all of the display cells in the panel so as to activate the individual display cells and to keep the display characteristics uniform, and is inserted at certain constant intervals. The all-erasure period AE is provided for erasing previous display contents before the address operation and sustaining operation for image display are newly commenced, through application of an erasure voltage pulse to all of the display cells on the panel.\nFIG. 14 shows an exemplary circuit diagram of the address drive circuit used for the address driver 202 of the plasma display device shown in FIG. 12. The address electrode drive circuit shown in FIG. 14 is disclosed typically in Patent Document 1 below. The address electrode can be assumed as a capacitive load 100 as shown in the drawing, because most of the current components flow through the address electrodes (A1 to Ad) are charge-discharge current for parasitic capacitance of the electrodes.\n[Patent Document 1] Japanese Patent Application Laid-Open No. 5-249916\nAssuming now that the number of the address electrodes (A1 to Ad) of the display panel reaches as much as 3072 (1024 pixels×RGB), the circuit is configured so that the outputs of 24 drive ICs, each having 128 address electrode drive circuit shown in FIG. 14 integrated therein, are connected to the address electrodes. The connection to the address electrodes is made typically through a flexible substrate on which the drive ICs are mounted in a unit of a single chip or two or more chips.\nIn the circuit diagram shown in FIG. 14, a single address electrode corresponds to the capacitive load 100. A low-side output element (N-channel MOSFET) 101 is connected between the ground, which is a low-voltage-side reference potential, and the capacitive load 100. A high-side output element (P-channel MOSFET) 102 is connected between a drive power source 107, capable of supplying high-voltage potential Va which corresponds to high-level of the address drive voltage, and the capacitive load 100.\nAs an exemplary circuit for driving the high-side output element 102, FIG. 14 shows a level shift circuit 108. The level shift circuit 108 drives the high-side output element 102 with the aid of an inverter circuit which comprises a P-channel MOSFET 103 and an N-channel MOSFET 104. The P-channel MOSFET 103 is driven by flip-flop operation of another inverter circuit which comprises a P-channel MOSFET 105 and an N-channel MOSFET 106.\nOperation in the address period ADD (FIG. 13) will be explained referring to a timing chart shown in FIG. 15. FIG. 15 shows timing relations between output voltage Vo of the address electrode drive circuit and input voltages VG1 to VG6 of the individual drive elements 101 to 106. In a rise-up period AT of output voltage Vo, the low-side output element 101 is cut off by inverting input voltage VG1 from high level to the low level, and the high-side output element 102 is turned on by inverting input voltage VG4 into the high level and input voltage VG6 into the low level. This adjusts output voltage Vo to high voltage potential Va. On the contrary, in the decay period TB of output voltage Vo, the high-side output element 102 is cut off by inverting input voltage VG4 into the low level and input voltage VG6 into the high level, and the low-side output element 101 is turned on by inverting input voltage VG1 into the high level. This adjusts output voltage Vo to the ground level (0V).\nThe rise-up time and decay time as seen in the waveform of output voltage Vo shown in FIG. 15 refer to times for charging and discharging the load capacitance 100 with the output currents from the high-side output element 102 and low-side output element 101. It is to be noted herein that the adjacent address electrodes on the display panel 201 shown in FIG. 11 are switched between high-voltage potential Va and the ground level depending on image to be displayed. The load capacitance 100 in this case therefore can largely vary and affect as follows. For the case where the adjacent electrodes placed on the left and right sides of a target address electrode are switched at the same time in the same direction (i.e., together from the ground level to high-voltage potential Va), the load capacitance 100 will have a minimum value as being contributed only by capacitance Cg without including capacitance Ca between the adjacent electrodes which are in no need of charge nor discharge. On the contrary, for the case where the left and right adjacent electrodes are switched at the same time in the opposite directions (i.e., one varies from the ground level to high-voltage potential Va, and the other varies from high-voltage potential Va to the ground level), the capacitive load 100 will have a maximum value of Cg+4Ca because both capacitances Ca between the adjacent electrodes will be supplied with doubled electric charge. Ratio of variation in the load capacitance 100 therefore reaches three times or more in general. Output currents from the high-side output element 102 and low-side output element 101 must be designed as being sufficiently large so as to successfully obtain a drive speed necessary for the display panel 201. This, however, results in a sharp change in the waveform of output voltage Vo under the minimum load and consequently in a drastic decrease in the transition time, and therefore raises a problem of increase in unnecessary electromagnetic radiation ascribable thereto. Interference of any other electronic instruments caused by unnecessary electromagnetic wave is known as EMI (electromagnetic interference), and is necessarily suppressed to a level allowable by any specified standards. In general procedures for suppressing unnecessary radiation, lack of suppressive measures in the initial stage of the design will unfortunately increase costs for any additional electromagnetic shield or filter elements."}
{"text": "This invention relates to flash memory storage, and more particularly to a serial-bus controller integrated with a parallel flash memory.\nFlash memory has gained wide acceptance for its non-volatile storage, which is ideal for portable devices that may lose power, since the data is not lost when stored in the flash memory. Flash memories are constructed from electrically-erasable programmable read-only memory (EEPROM) cells.\nRather than use a randomly-addressable scheme such as is common with dynamic-random-access memory (DRAM), many flash memories use a block-based addressing where a command and an address are sent over the data bus and then a block of data is read or written. Since the data bus is also used to send commands and addresses, fewer pins are needed on the flash-memory chip, reducing cost. Thus flash memory is often used as a mass-storage device rather than a randomly-addressable device.\nUniversal-Serial-Bus (USB) has become a popular standard interface for connecting peripherals to a host such as a personal computer (PC). USB-based flash-memory storage devices or “drives” have been developed to transport data from one host to another, replacing floppy disks. While large external flash drives may be used, smaller USB flash drives known as key-chain or key drives have been a rapidly growing market.\nA USB flash-memory device can be constructed from a microcontroller, a flash-memory controller or interface, and one or more flash-memory chips. A serial interface on the microcontroller connects to the USB bus to the host, and data from the serial interface is transferred through the microcontroller to then flash controller and the written to the flash-memory chips.\nThe microcontroller usually contains an internal ROM with a control program that is read by the internal central processing unit (CPU) of the microcontroller when the microcontroller is booted or powered up. Once initialized with the control program, the CPU can control data transfers between the serial interface and the flash controller.\nSometimes the user may desire to connect to more than one USB flash-memory device. The user can install a USB hub, and then plug the USB flash-memory devices into the USB hub's downstream ports. USB hubs allow one USB port on a host to fan out to multiple end USB devices or endpoints. A basic USB hub has a repeater that repeats data from the host to all down-stream devices, while more intelligent hubs based on the USB 2.0 standard can buffer data to different down-stream ports.\nFIG. 1 shows a prior-art USB hub that connects to multiple flash-memory USB endpoint devices. Host 10 includes USB host controller 12 that generates transactions to USB devices over USB bus 18 using the USB protocol. USB hub 20 is connected to a cable containing USB bus 18. USB hub 20 fans out USB bus 18 to several downstream USB devices that connect over additional USB bus segments.\nThree USB flash-memory systems 14, 15, 16 are connected to USB hub 20 by USB bus segments. USB flash-memory system 14 can be accessed by USB host controller 12 through USB hub 20. Since USB hub 20 passes all host transfers through to downstream devices, USB flash-memory system 15 is visible to host 10 as a second flash drive, while USB flash-memory system 16 is visible to the host as a third flash drive.\nUsers may be able to purchase a single USB flash-memory drive with a larger storage capacity and directly connect it to USB bus 18, without the need for USB hub 20. Some board manufacturers may integrate USB hub 20 together with USB flash-memory systems 14, 15, 16 on a single USB flash card or box. However, this can be expensive when USB flash-memory systems 14, 15, 16 are flash-memory chips, since each chip may have many pins. For example, a flash-memory chip with an 8-bit or 16-bit data bus may have 48 total pins. This can increase the size of the USB flash device. Power consumption is higher due to the large number of data lines in the parallel buses to each flash-memory chip.\nThe transfer bandwidth of each flash-memory chip is also somewhat limited. For example, a 50 nanosecond access time for an 8-bit flash chip has a bandwidth of 160 Mbps. While this is sufficient for the USB 1.1 standard that supports a data rate of 12 Mbps, the newer USB 2.0 standard supports a data rate of 480 Mbps. Thus a wider data bus, requiring more flash-memory chips, is needed to support the faster USB 2.0 speeds. However, the additional flash-memory chips and wider buses increases board size, manufacturing cost, and product size.\nWhat is desired is to integrate a microcontroller with a flash-memory array. It is desired to have a wide internal bus from the microcontroller to the flash-memory array to improve the data bandwidth while having few external pins to reduce cost and required board space.\nIt is further desired to eliminate the internal ROM on the microcontroller. Instead of booting from the internal ROM, it is desired to use a control program stored in the flash-memory array. However, it is also desired to use a block-addressed rather than a randomly-addressable array for the flash storage."}
{"text": "The present invention relates to a method and apparatus for final rethermalizing of precooked food products such as pasta products and particularly, individual sized portions of pasta products.\nThere is a need for an apparatus for rethermalizing food products such as pasta products quickly and to order, and particularly a need to rethermalize individual sized ready-to-eat portions of pasta promptly and to order. Even more particularly, there is a need to be able to supply individual sized portions of pasta which are cooked properly and which have the requisite texture and consistency on a \"fast food\" basis.\nIn the past, it has not been possible to prepare pasta quickly, with the required texture and in individual sized portions. Means are known for cooking large quantities of pasta, but these typically will result in cooked pasta products which do not have the proper texture when served. If pasta is cooked in large quantities, typically, after some time passes, the pasta will continue to cook and will become soft, and not have the desired texture and consistency. Other systems have been developed to provide individual sized portions of pasta, but these systems are typically cumbersome and complex and rely on high pressure and high temperature. If pasta is cooked according to the prior art systems on a made to order basis, the process typically takes too long, particularly for consumers who are \"on the run\" and need to order and obtain food on a speedy basis. Accordingly, there is a need in the food service industry for equipment which will allow a consumer to obtain an individual sized portion of a pasta product to order, cooked to the proper consistency and texture, and which rethermalizes the pasta product in a short period of time, i.e., in a time period of a little more than a minute and generally, less than about two minutes and no more than about 4-5 minutes.\nApplicant is aware of the following patents:\n______________________________________ Patent No. Dated Name ______________________________________ 5,347,917 9/20/94 Vezzani et al. 4,803,917 2/14/89 Barbieri 4,732,080 3/22/88 Vita 4,214,514 7/29/80 Contino et al. 5,215,001 6/1/93 Narcisi 5,033,364 7/23/91 Narcisi et al. 5,010,806 4/30/91 Narcisi et al. 4,869,160 9/36/89 Pratolongo 4,803,916 2/14/89 Tacconi 4,741,261 5/3/88 DiMaria 4,569,277 2/11/86 Stiglich 5,630,359 5/20/97 Stein 4,901,632 2/20/90 Lori 5,586,487 12/24/96 Marino 4,175,482 11/27/79 Kumagaya 3,790,391 2/5/74 Bolleter et al. 5,609,093 3/11/97 Hohler et al. 5,220,864 6/22/93 Ejiri et al. 5,313,876 5/24/94 Hilger et al. 4,821,632 4/18/89 Bolzani et al. 3,928,045 12/23/75 Tsunoda et al. 5,351,605 10/4/94 Sai et al. 5,241,899 9/7/93 Kuhlman 5,365,831 11/22/94 Kuhlman 5,613,423 3/25/97 Polster 5,282,498 1/1/94 Cahlander et al. 5,353,847 10/11/94 Cahlander et al. 5,315,919 5/31/94 Hoeberigs 5,404,797 4/11/95 Millar ______________________________________\nThe Vezzani et al. patent discloses a machine, which may be microprocessor controlled, for preparing pasta based dishes starting from pre-cooked deep frozen pasta products. The machine includes a basket adapted to receive dosages of pre-cooked deep frozen pasta, a vat containing heated water, a drive arrangement for moving the basket between the dipped position where the basket is submerged under the heating water contained in the vat for heating the pre-cooked deep frozen pasta and a removed position where the basket is removed from the vat.\nThe Vezzani et al. reference is directed in particular to the means employed for loading the pasta into the cooking basket. The means employed comprises two chests 12 each slidable in a respective box-type seat 13 inside a body 11. The chest is open at the top and bottom. The seats 13 have bottom openings 19 so that the pasta disposed in the chests 20 can drop into the basket 23 which can then be immersed in the heated water in the vat 20. A basket is immersed in the vat 20 using a pivoting arrangement.\nAnother reference is U.S. Pat. No. 4,803,917 to Barbieri for an automatic apparatus to cook and serve ready-to-eat pasta helpings. In this reference, the cooking vessel 5 is movable between a higher and lower position. In the higher position, the pasta contained in a basket 18 is cooked. The basket has a clam shell arrangement which opens from the bottom to dispense the pasta onto a plate, for example, plate 32 shown in FIG. 4, once the pasta has been cooked. This reference notes in col. 5, at lines 12-19 that it can be used for pre-cooked pasta, in which case the cooking time can be reduced. Another reference which discloses the cooking of pre-cooked foods including pasta is U.S. Pat. No. 4,732,080. This reference describes an apparatus for preparing hot food, including the pre-cooking operation and the later storing, heating and dispensing operation.\nThere are a number of references directed to cookers for cooking individual portions of food products, including pasta. For example, see U.S. Pat. No. 4,214,514 to Contino et al., U.S. Pat. Nos. 5,215,001, 5,033,364 and 5,010,806 to Narcisi and Narcisi et al., U.S. Pat. No. 4,869,160 to Pratolongo, U.S. Pat. No. 4,803,916 to Tacconi, U.S. Pat. No. 4,741,261 to DiMaria, U.S. Pat. No. 4,569,277 to Stiglich and U.S. Pat. No. 5,630,359 to Stein. Each of these references is directed to a pasta cooking apparatus. Generally, these devices employ pressurized heated water to cook the pasta quickly. Accordingly, uncooked pasta can be cooked in a relatively short period of time because it is cooked at a high temperature and pressure. There is a need for a less complex and less dangerous pasta rethermalizing device wherein pasta is rethermalized at atmospheric pressure and at temperatures below the boiling point of water for a short period of time because the pasta is precooked and therefore does not need to be cooked for a long period of time.\nOther references, not necessarily directed to cooking individual servings of pasta, but which employ pressurized cooking techniques include U.S. Pat. No. 4,901,632 to Lori, U.S. Pat. No. 5,586,487 to Marino, U.S. Pat. No. 4,175,482 to Kumagaya and U.S. Pat. No. 3,790,391 to Bolleter et al.\nIn addition, there are a number of other references of which applicant is aware including U.S. Pat. No. 5,609,093 to Hohler et al., U.S. Pat. No. 5,220,864 to Ejiri et al., U.S. Pat. No. 5,313,876 to Hilger et al, U.S. Pat. No. 4,821,632 to Bolzani et al., U.S. Pat. No. 3,928,045 to Tsunoda et al., U.S. Pat. No. 5,351,605 to Sai et al. and U.S. Pat. No. 5,241,899 and U.S. Pat. No. 5,365,831 to Kuhlman. These are directed to larger scale pasta cooking systems or pasta cooking systems using specialized techniques, for example, employable in vending machines. Other references include U.S. Pat. No. 5,613,423 to Polster for a cooker/rethermalizer, U.S. Pat. Nos. 5,282,498 and 5,353,847 to Cahlander et al. directed to food dispensing devices and U.S. Pat. No. 5,315,919 to Hoeberigs and U.S. Pat. No. 5,404,797 to Millar directed to apparatus for use with fried foods.\nAlthough the above patents show a number of devices and methods for cooking pasta and other food products, there is a need in the industry for a method and apparatus for cooking individual sized portions of pasta products, and particularly pre-cooked pasta products, and which can be used to provide ready-to-eat pasta products to consumers quickly and conveniently."}
{"text": "The Internet can provide a user with a wealth of information and computing power, but can also result in security vulnerabilities. One class of security vulnerabilities involves improper access to a particular URL (Uniform Resource Locator) security zone. A URL security zone includes a group of URL namespaces that are assigned an equal level of permissions (or trust). Each URL action for the zone has an appropriate URL policy assigned to it that reflects the level of trust given to the URL namespaces in that zone. The URL actions may include actions that a browser can take that might pose a security risk to the local computer. Actions that might pose a security risk include actions such as running a JAVA applet or an ACTIVEX control. Additionally, certain ACTIVEX controls determine that they are hosted in a security zone where possibly unsafe actions are permissible. A malicious script that gains access to this security zone may be able to repurpose the ACTIVEX control for impermissible actions. The URL policy that controls these actions within the security zone, determines what permission or trust level is set for a particular URL action. For example, the policy may dictate that the safety level for such ACTIVEX controls be set to high.\nCertain browsers divide URL namespaces into URL security zones, which are assigned different levels of trust. The default URL security zones may be customized by changing the URL policy settings for each URL action. Each URL security zone has a set of URL actions with a URL policy assigned to each action. A URL policy is assigned to each URL action to determine how that URL action is handled. A number of default security zones have been created to assist in the assignment of security policies to various browser actions, including the local intranet zone, trusted sites zone, Internet zone, restricted sites zone, and local machine zone.\nUsers use the local intranet zone for content located on an organization's intranet. Since the servers and information is within an organization's firewall, a user or organization can assign a higher trust level to the content on the intranet.\nUsers use the trusted sites zone for content located on Web sites that are considered more reputable or trustworthy than other sites on the Internet. Users can use this zone to assign a higher trust level to these sites to minimize the number of authentication requests. The user adds the URLs of these trusted Web sites to this zone.\nUsers use the Internet zone for Web sites on the Internet that do not belong to another zone. This default setting causes the Web browser to prompt the user whenever potentially unsafe content is ready to download. Web sites that are not mapped into other zones automatically fall into this zone.\nUsers use the restricted sites zone for Web sites that contain content that can cause, or may have previously caused, problems when downloaded. Users can use this zone to cause the Web browser to alert them whenever potentially unsafe content is about to download, or to prevent that content from downloading. The user adds the URLs of these mistrusted Web sites to this zone.\nThe local machine zone is an implicit zone for content that exists on the local computer. The content found on the user's computer, except for content that the Web browser caches on the local system, is treated with a high level of trust.\nAn entity attempting to improperly gain access to a user's computer, may attempt to get access through the local machine zone. Accordingly, the local machine zone presents a possible security vulnerability if unauthorized entities gain access to this security zone."}
{"text": "The goal of vegetable breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits may include any trait deemed beneficial by a grower and/or consumer, including greater yield, resistance to insects or disease, tolerance to environmental stress, and nutritional value.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant or plant variety. A plant cross-pollinates if pollen comes to it from a flower of a different plant variety.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants of different genotypes produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform varieties requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods that have been used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more plants or various other broad-based sources into breeding pools from which new lines and hybrids derived therefrom are developed by selfing and selection of desired phenotypes. The new lines and hybrids are evaluated to determine which of those have commercial potential."}
{"text": "In a vehicle having an automatic transmission, a clutch assembly smoothly engages a rotating engine crankshaft with a stationary driveshaft for transmission of power to the drive wheels, and also disengages the respective shafts to interrupt power transfer therebetween to permit, for instance, smooth shifting between various gears of a planetary gear set. Clutch assemblies or clutches are torque-transmitting devices typically having a series of friction elements, i.e. a clutch pack, located within a clutch housing and actuated by a clutch piston, the piston being powered or energized by a supply of hydraulic fluid. The hydraulic fluid supply is typically pressurized by a controllable pump. When hydraulic clutch pressure is reduced, the clutch is released or disengaged, and likewise, when clutch pressure is increased, the clutch is actuated or engaged. The hydraulic system may be further controllable to actuate other transmission components, such as a specific clutch or a series of clutches within an automatic transmission.\nOne such hydraulically-actuated clutch is a latched-pump applied clutch (LPAC). In an automatic transmission containing an LPAC, a controllable pump pressure engages the LPAC while a latching valve closed to thereby trap and substantially seal a supply of pressurized hydraulic fluid within the LPAC circuit. Once adequate clutch-apply pressure has been trapped or sealed within the LPAC circuit, controllable pressure in the main pump circuit may then be reduced as required in order to minimize spin losses elsewhere in the transmission without thereby diminishing available clutch-apply pressure within the LPAC circuit.\nA leak or series of leaks within a sealed LPAC circuit may cause a decrease or reduction in available clutch-apply pressure, the effect varying with the severity and/or number of leaks in the circuit. Leakage or bypass might occur at various points within the circuit, including around piston seals, valve body gaskets, various component connections, within latching valves, or through the inherent porosity of cast transmission components. Such leaks, particularly in a relatively stiff of low-compliance system, can deplete available clutch-apply pressure within the LPAC circuit. When applied to a rotating-type clutch in particular, that is a stationary piston is used to clamp or apply a rotating clutch pack, a plurality of thrust bearings are commonly used to allow relative motion without an undue increase in drag. Under these circumstances, however, thrust bearing spin loss may increase along with the increase in clutch-apply pressure."}
{"text": "(1) Field of the Invention\nThis invention relates to a roll-over suppressing control apparatus for a vehicle for suppressing overturning of the vehicle while the traveling performance of the vehicle is secured.\n(2) Description of the Related Art\nWhen a vehicle turns, depending upon the traveling speed or the steering situation of the vehicle, a moment in the rolling direction is generated on the vehicle and sometimes degrades the stability of the vehicle. In such an instance, if the moment in the rolling direction increases excessively, then there is the possibility that the vehicle may be placed into roll-over (overturning). Therefore, various control apparatus have been proposed conventionally which take a braking step before the vehicle is placed into roll-over.\nFor example, Japanese Patent Laid-Open No. 2000-104582 (hereinafter referred to as Patent Document 1) discloses a vehicle behavior control apparatus having such a configuration that, if it is decided that the vehicle behavior is unstable, then fuel supply to an engine is blocked and, after the vehicle behavior is placed back into a stable state, the fuel supply to the engine is re-started in response to the value of the lateral acceleration (lateral G) of the vehicle. According to the configuration, in a state wherein the vehicle behavior is unstable, the engine torque is reduced so that the wheels may grip the road surface well, and after the vehicle is placed back into a stable state, the control of reducing the engine torque is ended in response to the value of the lateral acceleration of the vehicle to restore a stable state wherein no hunting of the control occurs.\nWhile the technique disclosed in Patent Document 1 takes a braking step of reducing the propelling force of the vehicle through control of the engine torque to suppress the vehicle speed, also another technique of applying braking force to the wheels to brake the vehicle has been proposed.\nFor example, Japanese Patent Laid-Open No. Hei 10-24819 (hereinafter referred to as Patent Document 2) discloses a braking force control apparatus having such a configuration that it includes a first control mode wherein it controls the vehicle behavior based on a turning characteristic (magnitude of the yaw rate) of the vehicle and a second control mode wherein rolling (overturning) is controlled based on an excessive roll index (lightness or decreasing speed of the load to a wheel) of the vehicle. In the first control mode, a difference in braking force is provided between turning inner wheels and turning outer wheels to control the yaw moment of the vehicle in the turning direction or turning round direction, but in the second control mode, braking force is applied to the wheels and the fuel amount to be supplied to the engine is reduced to control the vehicle speed. According to the configuration, when the vehicle is in the first control mode, recovery of the vehicle behavior can be achieved, and when the vehicle is in the second control mode, the centrifugal force to act upon the center of gravity of the vehicle can be reduced to prevent rolling and the traveling stability can be maintained well.\nIn this manner, according to the apparatus disclosed in Patent Document 2, in order to prevent rolling of the vehicle, the control of applying braking force to the wheels is carried out simultaneously in addition to the control of decreasing the engine torque disclosed in Patent Document 1.\nIncidentally, turning of a vehicle is divided into non-steady (transient) turning upon lane change, emergency avoiding steering and so forth wherein the turning direction is changed over suddenly and steady turning represented by turning along a loop bridge or the like along which turning in the same direction continues.\nIn the non-steady turning, a high lateral acceleration does not appear immediately after a changeover of the turning direction, but a lateral acceleration appears rather later. However, a roll rate appears immediately after a changeover of the turning direction and generally is less liable to be controlled by the driver. If the roll rate becomes excessively high during such turning, then the vehicle may possibly turn over to its side even if the lateral acceleration remains low.\nOn the other hand, upon steady turning, a high roll rate does not occur with the vehicle except when the turning behavior of the vehicle changes such as upon starting of turning or upon ending of turning, and the vehicle can be controlled readily by the driver. However, a lateral acceleration in the turning outer side direction normally appears. Then, if the lateral acceleration becomes excessively high during such turning, then the vehicle may possibly turn over to its side even if the roll rate remains low.\nTaking notice of such different types of turning of a vehicle as described above, if the control of applying braking force to the wheels as in the technique disclosed in Patent Document 2 and the control of reducing the engine torque are carried out simultaneously, then upon non-steady turning such as upon emergency avoiding turning, the vehicle can be slowed down sufficiently to suppress the roll-over and can be turned safely. However, upon steady turning, since the control of the vehicle during turning by the driver is comparatively easy, the vehicle is sometimes slowed down excessively by the driver. Further, since the engine output power is suppressed, the intention of the driver to raise the vehicle speed or to keep the vehicle speed becomes less likely to be reflected on the vehicle speed, and there is a subject also that the driving feeling is deteriorated."}
{"text": "As a method for manufacturing a semiconductor device, a method has been conventionally known of sealing one or more semiconductor chips fixed to a substrate, etc. with a sealing resin and dicing the sealed body to form a packaged semiconductor device unit. For example, a sealing sheet constituted with a thermosetting resin has been known as the sealing resin (for example, refer to Patent Document 1)."}
{"text": "This invention relates to certain herbicidal ethers, agriculturally suitable compositions thereof, and a method for their use as broad spectrum preemergent or postemergent herbicides.\nNew compounds effective for controlling the growth of undesired vegetation are in constant demand. In the most common situation, such compounds are sought to selectively control the growth of weeds in useful crops such as cotton, rice, corn, wheat and soybeans, to name a few. Unchecked weed growth in such crops can cause significant losses, reducing profit to the farmer and increasing costs to the consumer. In other situations, herbicides are desired which will control all plant growth. Examples of areas in which complete control of all vegetation is desired are areas around railroad tracks, storage tanks and industrial storage areas. There are many products commercially available for these purposes, but the search continues for products which are more effective, less costly and environmentally safe.\nPayne et al., U.S. Pat. No. 4,567,283 and Payne et al., U.S. Pat. No. 4,670,041 disclose a variety of herbicidal bicyclic ethers of the Formula ##STR1##\nU.S. Pat. No. 4,798,621 and WO 8,902,219 both disclose bicyclic ethers and their method-of-use in rice.\nU.S. Pat. No. 4,486,219 discloses bicyclic ethers of the formula: ##STR2##\nWO91/03464 discloses bicyclic ethers such as: ##STR3## wherein inter alia X and Y are independently C(R.sup.3) (R.sup.4)OR.sup.5 ;\nR.sup.5 is C.sub.1 -C.sub.3 alkyl, C.sub.2 -C.sub.4 alkenyl, C3-C.sub.4 alkynyl, CH.sub.2 -phenyl or C.sub.1 -C.sub.4 alkyl substituted with OR.sup.8 or OCF3; PA1 R.sup.2 is H, C.sub.1 -C.sub.4 alkyl, C.sub.2 -C.sub.4 alkenyl or C.sub.2 -C.sub.4 alkynyl; PA1 R.sup.3 and R.sup.4 are independently H, C.sub.1 -C.sub.4 alkyl, C.sub.2 -C.sub.4 alkenyl, C.sub.2 -C.sub.4 alkynyl or C.sub.1 -C.sub.3 alkyl substituted with OCH.sub.3 or OCH.sub.2 CH.sub.3 ; ##STR5## X is CR.sup.5 R.sup.6 Y.sup.1, CR.sup.5 .dbd.CR.sup.6 Y.sup.3 or R.sup.5 R.sup.6 PA1 R.sup.5 and R.sup.6 are independently H or C.sub.1 -C.sub.3 alkyl PA1 Y.sup.1 is H; C.sub.1 -C.sub.4 alkyl optionally substituted with Y.sup.2 ; C.sub.3 -C.sub.6 cycloalkyl; CN; C(O)NR.sup.7 R.sup.8 ; C(O) NHOR.sup.9 ; CO.sub.2 R.sup.10 ; C(O)R.sup.11 C(OR.sup.22) (OR.sup.23)R.sup.11 ; CHO; CH(OR.sup.22) (OR.sup.23); CH-NOR.sup.16 CR.sup.11 .dbd.NOR.sup.17 ; SR.sup.15 ; halogen; C.sub.2 -C.sub.4 alkenyl optionally substituted with 1-3 halogens, CO.sub.2 (C.sub.1 -C.sub.3 alkyl), CN or CHO; C.sub.2 -C.sub.4 alkynyl optionally substituted with halogen; S(O).sub.n R.sup.12 ; SO.sub.2 NR.sup.13 R.sup.14 ; NO.sub.2 N.sub.3 SCN; phenyl optionally substituted with 1-3 substituents selected from CH.sub.3, OCH.sub.3, SCH.sub.3, CF.sub.3 and halogen; imidazole; triazole; tetrazole; or ##STR6## Y.sup.2 is CN; C(O)NR.sup.7 R.sup.8 ; C(O)NHOR.sup.9 ; CO.sub.2 R.sup.10 ; C(O)R.sup.11 ; CHO; CH.dbd.NOR.sup.16 ; CR.sup.11 .dbd.NOR.sup.17 ; SR.sup.15 ; 1-3 halogens C.sub.2 -C.sub.4 alkenyl optionally substituted with 1-3 halogens; C.sub.2 -C.sub.4 alkynyl optionally substituted with halogen; S(O).sub.n R.sup.12 ; SO.sub.2 NR.sup.13 R.sup.14 ; NO.sub.2 ; N.sub.3 ; SCN; phenyl optionally substituted with 1-3 substituents selected from CH.sub.3, OCH.sub.3, SCH.sub.3, CF.sub.3 and halogen; imidazole triazole tetrazole or ##STR7## Y.sup.3 is H, C.sub.1 -C.sub.3 alkyl, CN, CO.sub.2 R.sup.10, C(O)R.sup.11, CHO or halogen; PA1 R.sup.7 R.sup.8 R.sup.13 and R.sup.14 are independently H or C.sub.1 -C.sub.3 alkyl; R.sup.7 and R.sup.8 may be taken together to form a 5- or 6-membered ring and R.sup.13 and R.sup.14 may be taken together to form a 5- or 6-membered ring; PA1 R.sup.9, R.sup.12, R.sup.16 and R.sup.17 are independently C.sub.1 -C.sub.3 alkyl; PA1 R.sup.10 is H; C.sub.1 -C.sub.6 alkyl optionally substituted with 1-3 halogens, OR.sup.18, CN or phenyl optionally substituted with 1-3 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; C.sub.3 -C4alkenyl; or C.sub.3 -C.sub.4 alkynyl; PA1 R.sup.11 is C.sub.1 -C.sub.6 alkyl optionally substituted with 1-3 halogens, OR.sup.18, CN or phenyl optionally substituted with 1-3 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; C.sub.2 -C.sub.4 alkenyl; C.sub.2 -C.sub.4 alkynyl; or phenyl optionally substituted with 1-3 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; PA1 R.sup.15 is H, C.sub.1 -C.sub.3 alkyl or benzyl; PA1 n is 1 or 2; PA1 Q is CH.sub.2 W or ##STR8## q and r are independently 0-2; R.sup.18 is H or C.sub.1 -C.sub.3 alkyl; PA1 R.sup.19 is H, halogen, C.sub.1 -C.sub.3 alkyl, OR.sup.20, SR.sup.20 or CN; PA1 R.sup.20 is C.sub.1 -C.sub.3 alkyl or C.sub.1 -C.sub.3 haloalkyl; PA1 Z is CH.sub.2, NR.sup.21, O, S or may be CH and taken to form a double bond with an adjacent carbon; PA1 R.sup.21 is H or C.sub.1 -C.sub.3 alkyl; PA1 R.sup.22 and R.sup.23 are independently C.sub.1 -C.sub.3 alkyl or may be taken together to from a 5- or 6-membered ring; PA1 R.sup.24, R.sup.25 and R.sup.26 are independently H or C.sub.1 -C.sub.2 alkyl; PA1 W is phenyl optionally substituted with 1-3 substituents selected from halogen, C.sub.1 -C.sub.3 alkyl, C.sub.1 -C.sub.3 alkoxy, OH, CN, C.sub.1 -C.sub.3 haloalkyl, C.sub.1 -C.sub.3 haloalkoxy, C.sub.1 -C.sub.3 alkylthio, C.sub.2 -C.sub.4 alkenyl and C.sub.2 -C.sub.4 alkynyl or W is a 5-, 6- or 7-membered heterocyclic ring containing 1 to 3 heteroatoms selected from the group 0-2 nitrogens, 0-2 oxygens and 0-2 sulfurs, each ring optionally substituted with 1-2 substituents selected from halogen, CH.sub.3 and OCH.sub.3 ; PA1 1) in compounds of Formulas I and II, Y.sup.1 is other than H, C.sub.1 -C.sub.4 alkyl or halogen, and Y.sup.2 is other than halogen(s); PA1 2) in compounds of Formula IV, when Q is other than CH.sub.2 W, then Y.sup.1 is other than H or C.sub.1 -C.sub.4 alkyl; PA1 3) the sum of q and r is 0-2; and PA1 4) if the sum of q and r is O then Z is CH.sub.2. PA1 1. Compounds of Formulae I, II, III or IV wherein: PA1 3. Compounds of Preferred 2 wherein: PA1 4. Compounds of Preferred 3 wherein: PA1 6. Compound of Preferred 5 wherein: PA1 7. Compounds of Preferred 6 wherein the compound is of Formula I. PA1 8. Compounds of Preferred 6 wherein the compound is of Formula II. PA1 9. Compounds of Preferred 6 wherein the compound is of Formula III. PA1 10. Compounds of Preferred 6 wherein the compound is of Formula IV. PA1 (2-endo, 4-endo )-(.+-.)-4-[(2,6-difluorophenyl)-methoxy]-5-methyl -6-oxa-bicyclo[3.2.1 ]octane-2-acetonitrile; (2-endo, 4-endo)-(.+-.)-5-ethyl-4-[(2-fluorophenyl) -methoxy ]-6-oxabicyclo[3.2.1]octane-2-acetonitrile PA1 (2-endo, 4-endo )-(.+-.)-5-ethyl-4-(phenylmethoxy)-6-oxabicyclo [3.2.1 ]octane-2-acetonitrile (2-endo, 4-endo)-(.+-.)-4-[(2,6-difluorophenyl)-methoxy]-5-ethyl-6-oxabicyclo [3.2.1 ]octane-2-acetonitrile PA1 (2-endo, 4-endo)-(.+-.)-5-ethyl-4-[(2-fluorophenyl)-methoxy]-5-methyl-6-oxabicyclo[ 3.2.1]octane-2 -acetonitrile; PA1 (2-endo, 4-endo)-(.+-.)-4-[(2-fluorophenyl)methoxy]-5-methyl-6-oxabicyclo[3.2.1]oct ane-2-acetonitrile; PA1 (2-endo, 4-endo )-(.+-.)-2-(bromomethyl)-5-methyl-4-(phenylmethoxy)-6-oxabicyclo[3.2.1]oct ane; PA1 (2-endo, 4-endo)-(.+-.)-5-ethyl-4-[(2-fluorophenyl) -methoxy]-2-(iodomethyl)-6-oxabicyclo[3.2.1]-octane; PA1 (2-endo, 4-endo)-(.+-.)-4-[(2-chlorophenyl)methoxy]-5-methyl-6-oxabicyclo[3.2.1]oct ane-2-acetonitrile; PA1 (2-endo, 4-endo)-(.+-.)-2-(iodomethyl)-5-methyl-4-(phenylmethoxy)-6-oxabicyclo[3.2. 1]octane; and PA1 (2-endo, 4-endo)-(.+-.)-5-ethyl-2-[(methylthio)methyl]-4-(phenylmethoxy)-6-oxabicyc lo[3.2.1]octane.\nIt is an object of the present invention to provide compounds and compositions which exhibit herbicidal activity on a variety of weed species. It is a further object of the present invention to provide compounds and compositions that are herbicidally safe to rice, cereals and broadleaf crops. It is a feature of the present invention to furnish novel oxabicyclo ethers that exhibit useful herbicidal activity. These and other objects, features and advantages will become apparent with respect to the following description of the invention."}
{"text": "Typically, unloading stations for unloading rail cars are made to be in fixed locations. Such fixed unloading stations are often deep pits with concrete side walls with conveyers made to catch material dropped from the rail cars, conveying the material upward out of the pits to, for example, load into dump trucks for transporting elsewhere, to stock pile the material or to feed other conveyer systems. These unloading stations are very costly to build and maintain and, as a result, it is often difficult to justify expanding into new geographical areas to build new fixed unloading stations. Furthermore, due to the limited proximity locations of the unloading stations, their can also be high cost in transporting the bulk material in dump trucks from the unloading stations to a final destination.\nMany attempts have been made to correct the problem of unloading rail cars in locations without the use of fixed unloading stations. However, such attempts have not matched the efficiency of unloading rail cars at the fixed unloading stations. For example, U.S. Pat. No. 5,402,874, to Dahlin et al., discloses a conveyer coupled to a motor vehicle for conveying bulk material being dumped from the bottom of a rail car to another location, such as a dump truck. While in use, the conveyer remains coupled to the motor vehicle and is positioned laterally above the motor vehicle with a conveyer bottom end under the rail car and a conveyer upper end located above the dump truck. This conveyer system may be efficient for moving bulk material from a single rail car; however, time is wasted for re-set up between each rail car since the conveyer system requires retracting and re-set up to begin conveying the material from a subsequent rail car. Further, since the conveyer system is linear, time is wasted in positioning and aligning the conveyer bottom end to the appropriate location where the bulk material is dumped from the rail car. In addition, the linear type of conveyer system is limited to bottom dumping from a single shoot in the rail car and, further, the flow of materials from the shoot may have to be restricted in order to avoid over-flow of the material onto the linear type of conveyer system.\nU.S. Pat. No. 6,561,742, to Crawford et al., discloses a rail car unloading apparatus utilizing a trackhoe mounted to the top of a rail car. The trackhoe requires a ramp to drive onto the upper surface of the rail car and then is mounted to the upper surface with a support bracket assembly. The upper surface of the rail car is then opened and the trackhoe is used for removing the contents of the rail car. Although this system facilitates removing the bulk material at locations other than a fixed unloading station, this system is limited in the speed by which the trackhoe can remove the bulk material from the rail car. Further, this system is inefficient in initial set-up time and the set-up time required for unloading subsequent rail cars. Additionally, trackhoes cannot remove all the material from the rail cars and, therefore, additional steps often must be taken to completely clean the material from the rail car, causing additional man-power and wasted time.\nTherefore, it would be advantageous to develop an apparatus and method for unloading material from rail cars that is readily movable to different locations and is as efficient in unloading material from multiple rail cars as that which is employed for unloading rail cars at fixed unloading stations."}
{"text": "The reduction in memory cell size required for high density dynamic random access memories (DRAMs) results in a corresponding decrease in the area available for the storage node of the memory cell capacitor. Yet, design and operational parameters determine the minimum charge required for reliable operation of the memory cell despite decreasing cell area. Several techniques have been developed to increase the total charge capacity of the cell capacitor without significantly affecting the chip area. These include structures utilizing trench and stacked capacitors, as well as the utilization of new capacitor dielectric materials having higher dielectric constants.\nOne common material utilized for capacitor plates is conductively doped polysilicon. Such is utilized because of its compatibility with subsequent high temperature processing, good thermal expansion properties with SiO.sub.2, and its ability to be conformally deposited over widely varying topography.\nAs background, silicon occurs in crystalline and amorphous forms. Further, there are two basic types of crystalline silicon known as monocrystalline silicon and polycrystalline silicon. Polycrystalline silicon, polysilicon for short, is typically in situ or subsequently conductively doped to render the material conductive. Monocrystalline silicon is typically epitaxially grown from a silicon substrate. Silicon films deposited on dielectrics (such as SiO.sub.2 and Si.sub.3 N.sub.4) result in either an amorphous or polycrystalline phase. Specifically, it is generally known within the prior art that silicon deposited at wafer temperatures of less than approximately 580.degree. C. will result in an amorphous silicon layer, whereas silicon deposited at temperatures higher than about 580.degree. C. will result in a polycrystalline layer. The specific transition temperature depends on the source chemicals/precursors and the reactor used for the deposition.\nThe continued increase in circuit density continues to drive the thickness of the capacitor and gate layers in DRAM cells to smaller dimensions. Yet, the application of doped polysilicon films in DRAM memory cells demands high quality thin films. Historically, very thin polysilicon films (i.e., less than or equal to about 100 Angstroms) are fundamentally difficult to achieve, and the necessity of conductively doping such thin films is even more problematic. Generally, there are two basic approaches to doping thin film polysilicon, namely a) ion implantation, and b) in-situ doping during silicon film deposition. It has usually been easier to control the doping concentration through ion implantation than by in situ doping. However, ion implantation of dopants is typically limited to lower dopant concentrations, such as less than or equal to about 1.times.10.sup.16 atoms/cm.sup.3. In many applications, desired dopant concentrations to render polysilicon films suitably conductive exceeds 10.sup.20 atoms/cm.sup.3. Therefore, in situ doping (which is capable of much higher level doping without damage as compared to ion implantation) continues to be a desired form of processing for thin polysilicon films.\nIn-situ doping is commonly achieved by feeding a dopant gas to a chemical vapor deposition reactor simultaneously with feeding of a suitable silicon source gas. Example silicon source gases include silane (SiH.sub.4), disilane (Si.sub.2 H.sub.6), trisilane, and organosilane, while example conductivity enhancing dopant gases include phosphine (PH.sub.3), diborane (B.sub.2 H.sub.6), arsine (AsH.sub.3) and certain organometallic precursors. However, the resultant silicon film from the high level doping can suffer serious drawbacks. Such include the degradation of the interface of such film to the substrate, which in the presence of the dopant atoms can initiate local undesired strains in the resultant silicon film. Such local strains can further propagate to induce gathering of dopant atoms which can be particularly problematic the thinner the desired resultant polysilicon film. The situation is magnified where a polysilicon film is being provided over a SiO.sub.2 layer, which is typically the case. The result of the above mentioned situations can be a non-uniformity of the film thickness and in the sheet resistance from position to position and in wafer to wafer.\nAccordingly, needs remain for producing improved polysilicon layers which are effectively in situ conductively doped during provision of the silicon on the substrate. Although the invention principally arose out of concerns specific to the art area of provision of thin doped polysilicon films, the artisan will appreciate applicability of the invention to thicker polysilicon films as well. The invention is limited only by the accompanying claims appropriately interpreted in accordance with the Doctrine Of Equivalents."}
{"text": "The present disclosure relates to an electric storage unit group, a charger, an electronic device, an electric vehicle, a method for charging the electric storage unit group, a method for discharging the electric storage unit group, a method for supplying and receiving power, and a method for determining a charging/discharging route in the electric storage unit group.\nIn order to meet requirements for high-power, high-capacity devices, a case of utilizing a type of a battery pack (assembled battery) configured by connecting a plurality of rechargeable batteries (single cells) in parallel and connecting a plurality of parallel-connected rechargeable battery groups in series is increasing. Such a battery pack is known from, for example, Japanese Patent Application Laid-Open No. H6-283210. Typically, in such a battery pack, charging power is supplied from one end of the battery pack. Also, a load is connected to the battery pack, and power is supplied to the load."}
{"text": "A sugarcane harvesting machine of this type is known from U.S. Pat. No. 3,830,046. An inclined conveyor is situated at the end of the product flow through this machine, and conveys sugarcane sections, formed by comminuting the cut sugarcane stalks in the machine, into a collection car. The inclined conveyor has two continuous traction mechanisms, designed as chains, which extend in parallel to one another and are connected by steel carriers. The chains run over upper drive wheels, and lower deflection wheels which are carried along and situated in the feed area of the sugarcane sections. The inclined conveyor is designed as a scraper conveyor; i.e., it has a robust sheet steel floor which forms a running surface over the entire length of the inclined conveyor and scrapes the carriers, thus upwardly conveying the sugarcane sections resting on the sheet steel floor.\nA first and a second extractor which are used to remove waste and the like from the sugarcane sections are situated in the sugarcane harvesting machine. The extractors operate via a rotor which generates an air flow which carries along the waste and conveys it to the outside. A first extractor is situated in a transition zone of the sugarcane sections on the inclined conveyor, while the second extractor is situated at the end of the inclined conveyor and removes waste from the sugarcane sections when they drop down at the upper end of the inclined conveyor and fall into the collection car.\nSuch inclined conveyors described above with reference to the prior art have many drawbacks. The steel chains wear out very quickly due to adhering sand and dust. In addition, they fall out of the division after even a few hours of operation due to the acting tension, causing the teeth to come out of the drive wheels. For this reason, the conveyor requires continuous monitoring and maintenance, and in spite of this has a short service life. A high level of wear and high energy consumption occur during operation of the inclined conveyor, since the sugarcane sections are scraped during conveyance over the steel floor of the conveyor. This is also associated with a high noise level due to the metal-to-metal contact. Furthermore, the scraper floor of the conveyor, which is made of solid sheet steel, is very heavy. A further drawback is that, due to the scraper floor which acts as a barrier, the inclined conveyor, although it is perforated, conveys almost all of the impurities along with it, resulting in only a very low screening and cleaning effect and thus subjecting the second extractor to heavy load."}
{"text": "(1) Field of the Invention\nThe invention is in the field of tools which are used in the removal of airplane passenger seats and the like from seat tracks which are typically mounted in the floor of an airplane. The invention is also in the field of tools which are used in the installation of passenger seats and the like in the seat tracks.\n(2) Description of the Related Art\nPassenger seats and the like are held in place in many commercial airplanes by being locked into seat tracks which are typically mounted in the floor of the airplane. In many instances, such as when the airplane is being converted from a passenger-carrying aircraft into a freight-carrying aircraft, it is desired to remove some or all of the seats.\nThe present method of unlocking the passenger seats, practiced by airline companies throughout the world, is by prying up on the seat locking plug (which is part of the aft seat locking mechanism) using a conventional flat-blade screwdriver and a hammer. This crude method often damages the seat track and/or the aft seat locking mechanism. As a result of the damage, the seat may become unremovable, thus creating a significant problem for the airline."}
{"text": "The present invention relates to a process for purifying hepatitis B viral surface antigen containing a preS2 peptide; and, more specifically, to a process for purifying hepatitis B viral surface antigen containing a preS2 peptide from a recombinant yeast cell, which includes a cell disruption step carried out in the presence of a chaotropic salt to obtain a cell homogenate and the step of alkalifying the cell homogenate to enhance the solubilization of the surface antigen.\nHepatitis B is one of the worldwide public health problems and approximately 200 to 300 millions of the world population are said to carry hepatitis B virus(xe2x80x9cHBVxe2x80x9d). The HBV infection frequently progresses into cirrhosis and hepatocellular carcinoma, leading to the possible death of the patient.\nHitherto, no treating agent for hepatitis B has been developed, and, as such, the importance of vaccines has been emphasized.\nBlumberg et al. discovered the so-called Australian antigen from the blood of hepatitis B patients in 1955; and Krugman et al. reported, in 1971, an active immunization method using a heat-treated human serum containing HBV, thereby offering the possibility of developing hepatitis B vaccines. Thereafter, the first generation hepatitis B vaccines, which are prepared by separating and purifying a hepatitis B viral surface antigen(HBsAg) from the blood plasma of hepatitis B patients, have been commercialized(M. R. Hilleman et al., Develop. Bio. Standard, 54, 3-12(1983)).\nHowever, the vaccines derived from the blood plasma have the problems that: their purification and inactivation processes are cumbersome and require high costs; supply of human blood is limited; and an inoculated person may be infected with the pathogens from the blood source.\nAccordingly, in order to solve the above problems, genetic engineering approaches have been tried in developing hepatitis B vaccines.\nFor instance, Valenzuela et al. have developed a process for producing HBsAg in yeast(Nature, 293, 347-350(1982)). The recombinant HBsAg(r-HBsAg) consists primarily of S protein(P25) having 226 amino acids, and when purified, it forms surface antigen particles which are almost identical to those of HBsAg separated from blood plasma.\nK. H. Heermann et al. have reported that the hepatitis B viral envelope protein contains significant amounts of L-protein(preS1+preS2+S: p39) and M-protein(preS2+S: p31), as well as S-protein(J. Virol., Nov., 396-402 (1984)). The preS1 peptide has been known to play an important role with respect to the onslaught of hepatitis B virus on the liver after it infects human. The preS2 peptide, which consists of 55 amino acids, has been found to help the antibody formation against the surface antigen in animal experiments(D. R. Milich et al., Proc. Natl. Acad. Sci. U.S.A., 82, 8168-8172 (1985)).\nFurther, it has been known that antibodies formed against the preS2 peptide exhibit defensive activity against viral infection(Y. Ito et al., Proc. Natl. Acad. Sci. U.S.A., 83, 9174-9178(1986)). Therefore, a vaccine containing the preS2 peptide may be useful to a person who cannot form antibodies against a pre-existing surface antigen. The development of such a vaccine is also important for the protection against infection by recently discovered hepatitis B viral variants.\nHowever, since the pres peptide is very sensitive to proteases present in a yeast cell, preparation of a hepatitis B viral surface antigen containing the pres peptide has met with various difficulties. In order to overcome the difficulties, Kobayashi et al. have produced a vaccine which is prepared by genetically modifying the protease-sensitive region between the preS2 and S peptides(J. Bacteriology, 8, 1-22(1988)); and U.S. Pat. No. 4,742,158 discloses a process for producing a vaccine containing the pres peptide, wherein the peptide is protected from protease attack by using a protease inhibitor in the cell disruption step. However, the effect of the genetic modification by Kobayashi et al. on the activity of the preS2 peptide has not been fully characterized, and the latter process is not practical because protease inhibitors are too expensive to be used in a mass purification process. Further, the level of preS2 peptide cannot be maintained beyond a certain amount regardless of the amount of protease inhibitors added when the purification time becomes longer as the purification scale or requirement increases.\nEuropean Patent No. 0 337 492 A1 provides a process for purifying HBsAg from the culture of Pichia sp. using potassium thiocyanate. Potassium thiocyanate is used for raising the yield of lipophilic proteins, and the entire process is aimed at the purification of the HBsAg containing the S peptide only. When the HBsAg further contains the preS2 peptide consisting of 55 amino acids in front of the S peptide consisting of 226 amino acids, it has immunological properties similar to those of the surface antigen comprising the S peptide only because the antigenicity and immunogenicity of the S peptide moiety are the same. However, the two surface antigens are inevitably different in their physicochemical properties. In particular, the preS2 peptide is sufficiently hydrophilic to be exposed on the surface of the antigen particle, and, accordingly, its purification requires a special procedure.\nTherefore, various processes for purifying the hepatitis B viral surface antigen comprising the preS2 peptide have been developed.\nEuropean Patent No. 0 130 178 A1 describes a process for purifying HBsAg comprising the preS2 peptide, which is characterized by separating the surface antigen using two liquid phases which are prepared by adding a suitable amount of dextran and glycol to a yeast extract. However, this process has problems in that: the separated surface antigen is not sufficiently pure; it is not suitable for a large scale purification process because dextran and polyethylene glycol are expensive; and it is not economical due to the fact that an additional procedure is required for the removal of a surfactant used.\nU.S. Pat. No. 4,742,158 teaches a process for purifying HBsAg containing the preS2 peptide, which comprises: preparing a yeast extract in the presence of various protease inhibitors and purifying the surface antigen therefrom by a series of column chromatographic separation steps using an affinity column prepared by attaching a human serum albumin polymer to a gel matrix as well as a hydrophobic column eluting with a surfactant. However, this process has such difficiencies as: the protease inhibitors used in the process are very expensive; the affinity column is not suitable for mass purification; and a special procedure is required for the removal of a surfactant used in the hydrophobic column chromatography.\nM. Kobayashi et al. have also reported a process for purifying HBsAg comprising the preS2 peptide(J. of Biotechnology, 8, 1-22(1988). However, this process is not suitable for mass purification as it employs an immunoaffinity column which gives a low yield.\nKorean Patent No. 065305 presents a process for purifying HBsAg, which comprises a pH precipitation, and silica and anion exchange column chromatography. This process has the disadvantage that it is difficult to maintain the preS2 peptide content at a suitable level because the preS2 peptide is digested by proteases during the initial stage of the process.\nTherefore, there still exists a need for an efficient process for purifying HBsAg containing the preS2 peptide.\nAccordingly, it is a primary object of the present invention to provide a process for purifying hepatitis B viral surface antigen containing a preS2 peptide from a yeast cell in a sufficiently pure state to be directly incorporated into a vaccine.\nIn accordance with one aspect of the present invention, there is provided a process for purifying hepatitis B viral surface antigen containing the preS2 peptide, which is expressed in recombinant organism cells, characterized in that the recombinant organism cells are disrupted using a buffer containing a chaotropic salt to obtain a cell homogenate, followed by alkalifying the homogenate to a pH ranging from 11.0 to 13.5 in the presence of a surfactant to enhance the solubilization of the surface antigen."}
{"text": "Capacitance material level probes of the type to which the present invention is directed typically comprise a central rod surrounded by a first layer of insulation, with the rod tip being exposed to act as a measurement antenna when mounted within a storage vessel. A tubular guard shield surrounds the first insulation layer and is itself surrounded by a second layer of insulation. A portion of the tubular guard shield adjacent to the rod tip is exposed and is adapted to be electrically energized separately from the rod so as to direct probe energy outwardly into the tank rather than through any material which may be coated onto the probe body. The second insulation layer is clamped within a nipple which is adapted to be threaded into a gland on a material vessel wall with the probe rod and guard shield projecting into the tank as described.\nMeasurement probes of the described character heretofore proposed in the art typically include the insulation layers and the guard shield as separate tubular elements telescopically received over the central rod. To obtain desired sealing within a pressurized vessel, it has often been necessary to crimp or squeeze the tip of the tubular guard shield around the first insulation layer to close off a leakage path therebetween. The separate tubular insulation elements and assembly thereof to the probe rod are expensive, and it is now proposed to perform the assembly operation and provide the necessary insulation in a single step by injection molding insulation material surrounding and between fixtured guard shield and probe elements. However, considerable shrinkage is encountered with thermosetting plastics suitable for this application, and crimping of the guard shield tip, which is undesirable in any event because it requires an extra assembly operation, is not entirely effective.\nAn object of the present invention, therefore is to provide an improved capacitance-type probe of the character described which is adequately and securely sealed at the tip of the guard shield.\nAnother object of the invention is to provide an economical method of manufacturing such a probe."}
{"text": "1. Field of the Invention\nThis invention relates to belt buckles, and more particularly to seat belt buckles; but, as the name implies, is not limited to these uses and is applicable wherever a catch of the nature described below is needed.\nAn object of this invention is the provision of a catch which will allow instant escape from a vehicle if used as a seat belt buckle, or instant release of cargo if used as a carge buckle. When used as a seat belt, this catch can be adjusted to release with the slap of a hand or the push of an arm struggling to be freed. The advantages of this characteristic can be easily seen--in many accidents survival depends in escaping the vehicle instantly, i.e., if it is burning or about to fall from an escarpment. This adjustment is determined by the producer in advance of fabrication in the preferred embodiment, but the concept includes an adjustment in the produced article. By elevating the stud further from the bottom of the male member, more tilt is required to depress the dog and spring and thereby release the male member; also by shortening the stud the same objective is accomplished. Whether the mechanism is adjusted for hair trigger release or for deliberate release, it is a faster release than that available in conventional button release seat belts.\n2. Description of the Prior Art\nThe prior art is replete with buckles having housings, snaps, and levers of one type or another--i.e. C. J. Schumann's Snap Lock Buckle (U.S. Pat. No. 2,110,684)--and many safety belt buckles have been cited in the prosecution of the prior and pending application."}
{"text": "The present invention relates to a copier capable of specifying a minimum necessary area of a photoconductive area to be erased in relation to a paper size, document size, trimming size and magnification, so that latent images within the specified area only may be erased to produce a quality copy.\nIn a copier, it has been customary to reproduce a document by an electrophotographic process which includes steps of charging, exposing, erasing, developing and transferring. During the erasing step, eraser means of the copier performs front end erasure, side erasure and rear end erasure so as to remove needless ones of electrostatic latent images which are formed on a photoconductive element. A typical implementation for erasure control known in the art is optically erasing needless latent images on a photoconductive element based on the size of papers which is sensed by a paper size sensor, which is associated with a paper cassette. Another known implementation relies on a sensor sensitive to the size of a document.\nThe problem with any of such prior art erasure control schemes is as follows. Assuming that the erasure control is performed on, for example, a paper size basis, where the document size is smaller than the paper size, stains in those portions of a document feed belt, document pressing plate and others which are located outside of the document area are formed as images and transferred onto papers to degrade the reproduction. Meanwhile, where the document size is smaller than the paper size, even needless regions of the photoconductive element which will not be transferred to papers are developed by toner, aggravating toner consumption as well as the loads on a cleaning unit."}
{"text": "The invention relates to a method and a configuration for producing a sound.\nSuch a configuration and such a method are known from R. E. Donovan et al.: xe2x80x9cAutomatic Speech Synthesiser Parameter Estimation using HMMsxe2x80x9d, IEEE 1995, pages 640-643; hereinafter xe2x80x9cDonovan et al.xe2x80x9d This publication discloses the production of a decision tree, for its part permitting a cluster formation for the modeling of triphones. For this purpose, a series of questions, which relate directly to a phonetic context and can be answered with yes or no, is used. Each answer sets up a further subtree. Consequently, training data of naturally spoken speech are projected onto the branches and finally leaves of the decision tree.\nThe decision tree is used for calculating the leaves to be used in order to determine hidden Markov models for all possible triphones not covered by the training data.\nHidden Markov models (HMM) for the modeling of sounds are known from E. G. Schukat-Talamazzini: Automatische Spracherkennung xe2x80x94Grundlagen, statistische Modelle und effiziente Algorithmen [Automatic Speech Recognitionxe2x80x94Principles, Statistical Models And Efficient Algorithms], Vieweg and Sohn Verlagsgesellschaft mbH, Brunswick/Wiesbaden 1995, pages 125-139. In the linguistic production of a word, the constituent sounds are realized with variable duration and in varying spectral composition. Dependent on the rate and rhythm of the speech, each individual phonetic segment of the utterance is allotted an unpredictable number of feature vectors; each vector includes not only its phonetic content but also information components relating to the speaker, ambience and slurring, which make identification of the sounds significantly more difficult.\nThese conditions can be modeled in a simplified form by a two-stage process, as FIG. 1 shows by the example of the word xe2x80x9chabenxe2x80x9d [have]. For the phonemes of the word, a corresponding number of states 102 to 106 are reserved in the model and are run through along the direction of the arrow 101 for producing speech. At every time pulse, it is possible to remain in the current state or transfer to the succeeding state; the system behaves randomly and is determined by the transfer probabilities 107 to 111 depicted. For example, the state 103 belonging to the phoneme /a/ is adopted over a number of successive brief analysis intervals (on average over ten), whereas realizations of the plosive /b/ take less time.\nWhile the first stage described of the random process models the time distortion of different pronunciation variants, a second stage serves for sensing spectral variations. Linked to each state of the word model is a statistical output function that weights the phonetic realization alternatives. In the example of FIG. 1, for the production of the phoneme /a/ not only the actually matching phone class 113 but also the phone class 114 is permitted with a positive probability (here: 0.1). The phone class 118 is likewise permitted for the production of the phoneme /n/ with a probability of 0.3. The formalism described also allows a description of an optional sound elimination, expressed by the xe2x80x9cbridgingxe2x80x9d 119 of the state 105 by a direct transfer between the states 104 and 106. The bridging is given a probability of 0.2 by way of example.\nThe transfer probabilities of the hidden Markov model can be determined based on training data. When it has been fully trained, the HMM represents a blueprint for the production of sequences of sounds (cf. Schukat-Talamazzini, pages 127-139). One method of training the HMM is to use the Baum-Welch algorithm.\nHowever, in the method described in Donovan et al. it is disadvantageous that, in the decision tree, only the leaves respectively found are used for the sound modeling.\nIt is accordingly an object of the invention to provide a method and configuration for determining a representative sound, method for synthesizing speech, and method for speech processing that overcomes the hereinafore-mentioned disadvantages of the heretofore-known devices of this general type and that, when determining a representative sound (from a large number of sounds), accounts for not only a structure, devised according to predetermined criteria, but also a characteristic state criterion of this structure.\nWith the foregoing and other objects in view, there is provided, in accordance with the invention, a method for determining a representative sound based on a structure. The first step of the method is forming, from a sound, a structure having a characteristic state criterion. The next step is providing a set of sound models, each sound model having a representative with a plurality of quality criterion. The next step is determining, in the structure, a first sound model from the set of sound models matching a first quality criterion. The next step is determining a second sound model from the set of sound models depending on the characteristic state criterion of the structure. The next step is forming an overall quality criterion for each representative by assessing representatives of the first and the second sound model with regard to a second quality criterion in addition to the first quality criterion. The next step is determining a representative having an adequate overall quality criterion with regard to the first and second quality criteria as a representative sound.\nIn the method for determining a representative sound based on a structure that includes a set of sound models, each sound model has at least one representative for the modeled sound. In the structure, a first sound model, matching with regard to a first quality criterion, is determined from the set of sound models. Dependent on a characteristic state criterion of the structure, at least one second sound model is determined from the set of sound models. Representatives of the first sound model and of the at least one second sound model are assessed in addition to the first quality criterion with regard to a second quality criterion. From the representatives of the first and the at least one second sound model, that at least one representative which has an adequate overall quality criterion with regard to the first and second quality criteria is determined as a representative sound.\nIn accordance with how the structure is ordered, a search is conducted within the structure for a matching sound model for the sound to be produced. In this case, xe2x80x9cmatchingxe2x80x9d applies with regard to the first quality criterion, which is predetermined in particular by the structure.\nThe structure may be configured as a tree structure, preferably as a binary tree. Such a tree structure has nodes (for the embodiment of the sound models), branches (for the hierarchical subdivision of the sound models dependent on the criteria on the basis of which the structure is constructed) and leaves (nodes from which no further branch extends).\nThe structure constructed based on predetermined criteria is then used in order to determine, depending on the characteristic state of the structure (in particular the tree structure) at least one second sound model from the set of sound models. In this case, the characteristic state criterion in the structure may be a measure of distance from the first sound model. In the case of the binary tree as the structure, all the sound models within a predetermined distance from the first sound model may be regarded as second sound models. Here, the term xe2x80x9cdistancexe2x80x9d is not necessarily to be interpreted in the local sense; rather, the xe2x80x9cdistancexe2x80x9d may also concern a dimension of distance with regard to one or more predetermined criteria.\nWith the first sound model and a set of second sound models that satisfy the characteristic state criterion and have the predetermined distance from the first sound model, a second quality criterion is determined for the representatives of the sound models. The overall quality criterion for each representative is made up of the first and the at least one second quality criteria. A candidate among the representatives that is suitable for producing the sound is determined by its xe2x80x9cadequatexe2x80x9d overall quality criterion that is an overall quality criterion that has a predetermined minimum value.\nA development of this is that the overall quality criterion is an arithmetic combination of the quality criteria.\nIn particular, it is a development that a representative of the at least one second sound model is additionally provided with a reduced quality, resulting from the characteristic state criterion, in comparison with a representative of the first sound model. Here, the distance mentioned above of the at least one second sound model from the first sound model is preferably used: a great distance indicates a great deviation from the first sound model, originally found as matching according to the structure. For such a great distance, a kind of xe2x80x9cpenaltyxe2x80x9d is awarded for the representative of the at least one second sound model, this penalty having an effect on the overall quality criterion.\nOne refinement is that a best representative, with regard to the overall quality criterion, is determined as the representative sound.\nAnother refinement is that the structure is devised based on at least one of the following criteria:\na) linguistic criteria, for example grammatical properties of a word;\nb) articulatory criteria, in particular the location of the articulation; and\nc) phonetic criteria, for example classified on the basis of plosive or nasal sounds.\nAn additional development is that the sound may be formed as one of the following possibilities:\na) phoneme,\nb) biphone,\nc) triphone,\nd) syllable,\ne) combination of syllables,\nf) word,\ng) combination of words.\nIt is also a refinement that, whenever the sound represents a combination of sounds, a construction of the structure takes place on the basis of context criteria of the sound within the combination of sounds. This means in particular that, in the case of a combination of sounds, questioning of the kind\nxe2x80x9cIs the sound to the left of the current sound a vowel?xe2x80x9d\nor\nxe2x80x9cIs the sound to the left of the current sound a plosive?xe2x80x9d\ncan in each case be answered with YES or NO and consequently the context to the right or left is used for constructing the structure, in particular the binary tree.\nWithin an additional refinement, the sound model takes the form of a hidden Markov model (HMM).\nA development is that a selection of n representatives which, based on their overall quality criterion, are in each case best suited for determining the representative sound is represented. The number n can in this case be predetermined application-dependently; the best representatives are preferably presented in the form of a list.\nFurthermore, it is a refinement that the first sound model is determined for the determination of the representative sound as soon as no improvement of the first quality criterion can be achieved any longer by a further search in the structurexe2x80x94in particular the tree structure.\nAnother refinement is that the first sound model is determined as soon as fewer representatives than are required for the determination of the representative sound are available with respect to a sound model from the set of sound models with a possibly better first quality criterion. In other words, in the search for an adequately good first sound model, only those sound models that comprise more representatives than is predetermined by a threshold value are taken into account.\nWithin a further refinement, a fuzzy logic assessment is performed for each second quality criterion. As a result, it is ensured that for each second quality criterion an own assessment is performed, carried out in particular in a normalized manner with regard to comparability with other assessments. Fuzzy logic provides a suitable method for such a normalized assessment.\nIn accordance with a further object of the invention, the method described above can be used in speech processing, in particular in speech synthesis.\nAlso provided for achieving the object is a configuration for determining a representative sound based on a structure that comprises a set of sound models. Each sound model has at least one representative for the modeled sound. The configuration has a processor unit that is set up in such a way that\na) in the structure, a first sound model, matching with regard to a first quality criterion, can be determined from the set of sound models;\nb) dependent on a characteristic state criterion of the structure, at least one second sound model can be determined from the set of sound models;\nc) representatives of the first and the at least one second sound model can be assessed in addition to the first quality criterion with regard to a second quality criterion;\nd) that at least one representative which has an adequate overall quality criterion with regard to the first and second quality criteria can be determined as a representative sound.\nThis configuration is suitable in particular for carrying out the method according to the invention or one of its developments explained above.\nOther features that are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as embodied in a method and configuration for determining a representative sound, method for synthesizing speech, and method for speech processing, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawings."}
{"text": "1. Field of the Invention\nThis invention relates to a method of power generation, and, in particular, to a method of geothermal power generation.\n2. Brief Statement of the Prior Art\nThe successful geothermal power projects to date have exploited subterranean steam reservoirs. The more common subterranean brine reservoirs have not been developed, primarily because of the difficulties experienced in producing hot geothermal brine. The brine is difficult to produce because it forms scale deposits, chiefly calcium carbonate, when it is depressured. Attempts to avoid the scaling have included the use of down hole heat exchangers, however, the corrosive nature of the brine makes this approach impractical.\nIn my prior patent, U.S. Pat. No. 4,059,156, I disclose and claim a method for producing brine by injecting, into a column of brine, a low density lift medium which is sufficient to reduce the density of the resulting mixture, permitting it to be produced by the hydrostatic head of the brine column without any significant depressuring of the brine. Thereafter, the brine can be depressured above surface in equipment adapted for handling of slurries and scaling liquids.\nWhile the method described and claimed in my aforesaid patent is a substantial improvement in geothermal brine production, improvements are desirable in increasing the thermal efficiency of the process by providing increased recovery of energy from the produced brine."}
{"text": "Currently, in the Third Generation Partnership Project (3GPP), standardization of a 4G wireless communication system (LTE-Advanced) of a next generation type is in progress. In the 4G, it is under consideration to introduce techniques such as multi user-MIMO (MU-MIMO) and coordinated multi points transmission (CoMP) in order to improve a maximum communication rate and the quality at a cell edge.\nIn the LTE, a radio frame of 10 ms and 10 sub frames configuring the radio frame of 10 ms are defined. Each sub frame includes a control area used to transmit a control signal (PDCCH) and a data area used to transmit user data (PDSCH). For example, Patent Document 1 discloses this frame structure.\nMeanwhile, in the 3GPP, machine type communications (MTC) have been discussed. Generally, the MTC has the same meaning as machine to machine (M2M) and refers to communication between machines, which is not used directly by human. Usually, the MTC is performed between a server and an MTC terminal which is not used directly by human. In the MTC, communication may be performed at a long period such as one week or one month."}
{"text": "Field of the Invention\nThe present invention relates to an ink, an ink cartridge, and an ink jet recording method.\nDescription of the Related Art\nThe recording speed of ink jet recording apparatuses has been dramatically improved as a result of the technical advance. Such ink jet recording apparatuses have been gradually introduced into the field of business document printing in place of laser recording apparatuses. Ink jet recording apparatuses typically use, as inks for texts to be used for recording images on recording media such as plain paper, aqueous inks containing self-dispersible pigments as a coloring material in consideration of high optical density.\nIn order to improve the optical density of images to be recorded and image qualities such as character quality, various components have been studied for the aqueous inks containing self-dispersible pigments as a coloring material. For example, an ink enabling recording of high quality images has been disclosed. The ink includes a water-soluble organic solvent that is unlikely to stably disperse a self-dispersible pigment to destabilize the dispersion state of the pigment after evaporation of liquid components (Japanese Patent Application Laid-Open No. 2006-045514). Another ink enabling recording of high quality images has been disclosed. The ink includes a monovalent cation salt to destabilize the dispersion state of a pigment after evaporation of liquid components (Japanese Patent Application Laid-Open No. 2000-198955).\nHowever, the results of studies by the inventors of the present invention have shown that some aqueous inks containing a self-dispersible pigment as a coloring material cannot maintain the storage stability. In particular, as the application field of the ink jet recording method has been rapidly spread, inks are required to have much higher quality and also to have higher storage stability as compared with conventional inks.\nAn object of the present invention is thus to provide an aqueous ink containing a self-dispersible pigment as a coloring material but having markedly higher storage stability as compared with conventional inks. The present invention is also directed to provide an ink cartridge including the aqueous ink and an ink jet recording method using the aqueous ink."}
{"text": "Computer systems, such as hardware systems and software systems that run on computers often have undetected flaws that can be exploited by hardware attacks or software attacks, such as malicious computer programs that are received over the Internet or other communication networks. The hardware attacks and software attacks can include Trojans, viruses, worms, spyware, and other malware. Many existing computer security systems combat hardware attacks and software attacks by attempting to prevent the attacks from compromising any part of the computer system.\nProcessors in computer systems are designed to protect sensitive data in memory from both hardware attacks and software attacks. Some processors provide cryptographic mechanisms for encryption, integrity, and replay protection. Memory encryption protects the confidentiality of memory-resident data. Integrity protection prevents an attacker from causing any hidden modifications to the ciphertext (i.e., encrypted data, as opposed to plaintext which is unencrypted data) in memory. Replay protection eliminates any undetected temporal substitution of the ciphertext. In the absence of encryption, integrity, and replay protections, an attacker with physical access to the system can record snapshots of data lines and replay the data lines at a later point in time to modify the data lines and attack the computer system."}
{"text": "This invention relates to a quantitative test for the determination of mercaptan-based odorants in natural, synthetic, L.P. gases and the like.\nConventionally, mercaptan-based odorants are used in gas so that escaping gas can be recognized by members of the general public. These odorants are usually injected into the gas supply prior to the supply entering the distribution facilities of a town or the like and if too much odorant is present, then the slightest leak causes widespread odors thus resulting in unnecessary complaints or reports from the general public. It should be understood mercaptan-based odorants normally used are extremely penetrating in their effect and that it is necessary to maintain the levels within fairly strict limits.\nAt the present time, there are methods of attempting to assess the amount of mercaptan-based odorant present and the most common method is for the operator to smell the gas. This method is qualitative in nature and depends on the olfactory nerve. It therefore has large variations in readings, depending on the operator, and apart from this inherent error, the operator with a cold is not able to perform the test.\nOdorant measuring tubes have been used and this method involves the injection of a known volume of gas by a pump into a tube filled with a reactant which is absorbed on a granular material. The length of the developed stain due to the reaction is a measure of the odorant level, but in actual tests inconsistencies have been observed in the length of the stain which varies with the number of pump strokes and the concentration of odorant. Difficulty has also been observed in measurement of the actual length of the stain as there is no definite demarcation line between the end of the stain and the unstained material.\nTitration also gives good quantitative results, but a minimum time of 30 minutes is required to conduct a single test. Furthermore, the apparatus itself is elaborate and requires a relatively high degree of operator skill. It also requires gas pressure greater than those on normal domestic service as the gas has to pass through several fritted glass bubblers.\nLaboratory apparatus such as a gas chromatograph can of course be used with good quantitative results, but the size and high cost of the apparatus prohibit its use as a field instrument."}
{"text": "Today, as the computer programming languages, for example High Level Language (HLL), continue to grow rapidly, programmers often need to translate a computer program written in one HLL to another HLL. This translation may be necessary because of: lack of skilled programmers in the another HLL; when a user is not accustomed with concepts of the HLL in which the computer program is available; business decisions to switch underlying platforms; the desire of the programmers or the users to convert the computer program to another language; or in a scenario, where the computer may not be able to execute and support the available computer program, due to specific hardware or software requirements. Thus, in such scenarios the computer program in one HLL can be translated to another HLL.\nThe translation can be done manually by the programmers, but the manual translation is a laborious, time consuming, and an expensive process which is also prone to errors. To address this problem, automatic translators have been developed to translate programs written in one HLL (source program) to another HLL (target program).\nGenerally, for translating the source program to the target program using automatic translators, the source program is first given as an input to a parser. The parser checks for errors in the source program and further builds a data structure of the source program. Examples of errors include, but are not limited to, parsing errors, syntax errors, and compilation errors.\nThe parser during the parsing process throws a parsing exception and halts the parsing process, if any parser errors are identified while parsing a statement in the source program. The parsed source program is then given as input to the automatic machine translator (hereinafter referred as translator). The translator translates the parsed source program to the target program. The translator can throw a translation exception and interrupt the translation process if any translation errors are identified during translation. Hence in the above translation method, the parsing process and the translation process is interrupted and halted when a parser error or a translation error occurs.\nMoreover, as the translation is a machine translation it is also prone to errors. To identify such errors, the target program can be analyzed manually or automatically after the translation process. The errors in the target program can be easily identified while analyzing, however it is difficult for the user to identify the actual origin/source of the errors. Hence, if the error has occurred due to the error present in the source program then in such a scenario, the user have to identify the errors in the source program by manually analyzing the source program and the target program simultaneously. Thus, manually analyzing the source code and the target code is a time consuming and tedious job for the user.\nHence there exists a need to efficiently manage the transformation of computer programs."}
{"text": "1. Technical Field\nThe present invention relates generally to semiconductor technology and more specifically to reducing the number of steps in forming MirrorBit(copyright) Flash memory.\n2. Background Art\nDifferent types of memories have been developed in the past as electronic memory media for computers and similar systems. Such memories include electrically erasable programmable read only memory (EEPROM) and electrically programmable read only memory (EPROM). Each type of memory had advantages and disadvantages. EEPROM can be easily erased without extra exterior equipment but with reduced data storage density, lower speed, and higher cost. EPROM, in contrast, is less expensive and has greater density but lack erasability.\nA newer type of memory called xe2x80x9cFlashxe2x80x9d EEPROM, or Flash memory, has become extremely popular because it combines the advantages of the high density and low cost of EPROM with the electrical erasability of EEPROM. Flash memory can be rewritten and can hold its contents without power. It is used in many portable electronic products, such as cell phone, portable computers, voice recorders, etc. as well as in many larger electronic systems, such as cars, planes, industrial control systems, etc.\nIn Flash memory, bits of information are programmed individually as in the older types of memory, such as dynamic random access memory (DRAM) and static random access memory (SRAM) memory chips. However, in DRAMs and SRAMs where individual bits can be erased one at a time, Flash memory must currently be erased in fixed multi-bit blocks or sectors.\nConventionally, Flash memory is constructed of many Flash memory cells where a single bit is stored in each memory cell and the cells are programmed by hot electron injection and erased by Fowler-Nordheim tunneling. However, increased market demand has driven the development of Flash memory cells to increase both the speed and the density. Newer Flash memory cells have been developed that allow more than a single bit to be stored in each cell.\nOne memory cell structure involves the storage of more than one level of charge to be stored in a memory cell with each level representative of a bit. This structure is referred to as a multi-level storage (MLS) architecture. Unfortunately, this structure inherently requires a great deal of precision in both programming and reading the differences in the levels to be able to distinguish the bits. If a memory cell using the MLS architecture is overcharged, even by a small amount, the only way to correct the bit error would be to erase the memory cell and totally reprogram the memory cell. The need in the MLS architecture to precisely control the amount of charge in a memory cell while programming also makes the technology slower and the data less reliable. It also takes longer to access or xe2x80x9creadxe2x80x9d precise amounts of charge. Thus, both speed and reliability are sacrificed in order to improve memory cell density.\nAn even newer technology allowing multiple bits to be stored in a single cell is known as xe2x80x9cMirrorBit(copyright)xe2x80x9d Flash memory has been developed. In this technology, a memory cell is essentially split into two identical (mirrored) parts, each of which is formulated for storing one of two independent bits. Each MirrorBit Flash memory cell, like a traditional Flash cell, has a gate with a source and a drain. However, unlike a traditional Flash cell in which the source is always connected to an electrical source and the drain is always connected to an electrical drain, each MirrorBit Flash memory cell can have the connections of the source and drain reversed during operation to permit the storing of two bits.\nThe MirrorBit Flash memory cell has a semiconductor substrate with implanted conductive bitlines. A multilayer storage layer, referred to as a xe2x80x9ccharge-trapping dielectric layerxe2x80x9d, is formed over the semiconductor substrate. The charge-trapping dielectric layer can generally be composed of three separate layers: a first insulating layer, a charge-trapping layer, and a second insulating layer. Wordlines are formed over the charge-trapping dielectric layer perpendicular to the bitlines. Programming circuitry controls two bits per cell by applying a signal to the wordline, which acts as a control gate, and changing bitline connections such that one bit is stored by source and drain being connected in one arrangement and a complementary bit is stored by the source and drain being interchanged in another arrangement.\nProgramming of the cell is accomplished in one direction and reading is accomplished in a direction opposite that in which it is programmed.\nA major problem has been that the diffusion of the bitlines limits how closely the bitlines can be placed and this in turn limits how far the memory cells can be reduced in size.\nWith the urgency of reducing device size, a solution to this problem has been long sought but has long eluded those skilled in the art.\nThe present invention provides a method of manufacturing an integrated circuit, which includes providing a semiconductor substrate and depositing a charge-trapping dielectric layer and a gate dielectric layer over the semiconductor substrate. Bitlines are implanted closely in the semiconductor substrate and annealed using a rapid thermal anneal, which allows the bitlines to be more closely spaced than possible with the prior art. Wordlines and gates are formed and source/drain junctions are implanted in the semiconductor substrate. An interlayer dielectric layer is deposited and the integrated circuit completed."}
{"text": "This invention relates generally to distress signal detection systems and more particularly to systems of such type which are adapted to provide unattended monitoring of the 500 KHz international radio-telegraph distress frequency.\nAs is known, the International Convention For Safety of Life At Sea (SOLAS) requires that certain classes of vessels carry radio-telegraph auto alarm equipment for enabling the unattended monitoring of a distress signal transmitted by a distant ship. Such auto alarm equipment is designed to monitor the 500 KHz international radio-telegraph distress frequency and provide an audible warning signal at a location on the vessel where a watch-stander is generally present when a properly coded distress signal has been detected. Ideally, the distress signal consists of four-second dashes with 1-second spaced intervals. The auto alarm equipment should respond to detection of four of such dashes. After transmission of the coded distress signal the distant ship transmits the international distress SOS and then its position.\nIn order to allow for human inaccuracies in transmitting a properly coded signal and because atmospheric noise having a frequency component of 500 KHz will also be detected by the auto alarm equipment, such equipment must be capable of responding to dashes which are somewhat less or somewhat greater than four seconds in duration and to spaces which may be correspondingly greater than or less than one second. In particular, SOLAS requires that the auto alarm equipment be activated when dashes vary in length from 3.5 to 6 seconds and spaces vary in duration from 10 milliseconds to 1.5 seconds."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for preparing alkanolamines such as ethanolamines and isopropanolamines. More specifically, it relates to a process for particularly and selectively preparing monoalkanolamines.\n2. Description of the Related Art\nAlkanolamines can easily be obtained by reacting an alkylene oxide with aqueous ammonia, and also in an industrial scale, they have been manufactured by this process.\nHowever, the alkanolamines obtained by the above-mentioned process are present in the state of a mixture of a monoalkanolamine, a dialkanolamine and a trialkanolamine, and it is an important theme to control the production ratio of these components. Particularly in recent years, the demand for the monoalkanolamine is larger as compared with that of the dialkanolamine and the trialkanolamine, and a technique for preparing the monoalkanolamine in a high production ratio has been required.\nIn general, with regard to a method for manufacturing the monoalkanolamine in a high production ratio by, for example, the reaction of ethylene oxide and ammonia, the following matters are known. That is to say, the reactivity between ethylene oxide and ammonia is lower as compared with the reactivity between ethylene oxide and monoethanolamine or diethanolamine, and therefore the ratio of the reaction products depends upon the ratio of ammonia to ethylene oxide. Thus, as ammonia is used in large excess, the production ratio of the monoethanolamine increases (K. Weissermel and H. J. Arpe, Translation supervised by Mitsuaki Mukouyama, \"Industrial Organic Chemistry--Main Materials and Intermediates--\", Tokyo Kagaku Dojin, p. 149 (1978).\nHowever, if aqueous ammonia is used and a molar ratio of ammonia/the alkylene oxide is high, the production ratio of the monoalkanolamine can be increased, but the volumetric efficiency of a reactor deteriorates, and excess aqueous ammonia is required to be collected and recycled. In consequence, the energy unit requirement decreases and the load of the ammonia collection system and the water collection system increases inconveniently.\nIn order to solve such problems, there is a method in which the water content of the aqueous ammonia is minimized, but since the reaction of the alkylene oxide and ammonia proceeds by virtue of the function of water as a catalyst, such a method causes the deterioration of activity. Thus, as some measures to the above-mentioned inconvenience, there have been suggested a method of using a solid acid catalyst and increasing the reaction temperature (e.g., Japanese Patent Application Laid-open No. 47728/1974; Zh. Prikl. Khim., Vol. 56, p. 1966 (1983); and U.S. Pat. No. 4,438,281) and another method in which the reaction is carried out in a supercritical state (Japanese Patent Application Laid-open Nos. 13751/1984 and 33247/1984), but these methods have a problem that reaction pressure is high and so a high-pressure reactor is required, and another problem that the load of the ammonia collection system increases.\nOn the other hand, there are known methods in which when ammonia reacts with an alkylene oxide to prepare alkanolamines, a carbonate of ammonia is used, whereby the production ratio of the monoalkanolamine can be increased under relatively gentle conditions (British Patent No. 497,093 and U.S. Pat. No. 2,186,392).\nThe present inventors have intensively investigated a process in which ammonia reacts with an alkylene oxide in the presence of a carbonate of ammonia in accordance with the above-mentioned U.S. Pat. No. 2,186,392 to prepare the alkanolamines, and it has been confirmed that even if ammonia is not used in large excess relative to the alkylene oxide, the production ratio of the monoalkanolamine can be remarkably increased.\nNext, it was tried to separate, from a reaction solution, the monoalkanolamine, the dialkanolamine and the trialkanolamine which were products. That is to say, after ammonia, most of the carbon dioxide and water had been distilled off, the alkanolamines which were the products were separated from the resultant bottoms, and at this time, there was elucidated a problem, which was heretofore unknown, which impurities such as nitrogen-containing compounds such as carbamic acid and carbamic acid esters and nitrogen-containing heterocyclic compounds such as oxazolidone and N-hydroxyethylpiperazine were present in the alkanolamines."}
{"text": "A system-on-chip (SOC) is a type of integrated circuit (IC) that includes a processor and one or more other circuits. The SOC may be built using a single chip substrate. The processor executes program code and is coupled to the circuits so that the processor and the circuits operate cooperatively with one another. One example of an SOC is a chip level implementation of a computer or other data processing system. In some SOCs, the processor is coupled to programmable circuitry that may be used to implement various circuits.\nCreating a custom SOC requires expertise in both software development and hardware design. Traditionally, software development and hardware design have been separate fields of endeavor with each requiring specialized training and knowledge. In the case of an SOC that includes programmable circuitry, specialized knowledge of the particular architecture of the target IC has also been required. In consequence, the creation of a customized SOC has been a complex and time consuming undertaking requiring skilled software developers and skilled hardware designers working collaboratively."}
{"text": "Shopping carts and other types of carts have been constructed in a variety of ways. Most known carts have a wheeled frame and a basket connected to the frame.\nCarts have been known to have various disadvantages depending on their construction and production method. Some disadvantages that are prevalent in the field are inefficient use of plastic, metal and/or space; elevated pressure points on the basket or frame that may accelerate damage or wear; connection points that increase the chance of accumulation of food, liquids or other substances, and thus increase the chance of corrosion at the connection point while being generally unhygienic and difficult to clean; costly and/or lengthy manufacturing and assembling of the cart subcomponents; disadvantageous rigidity or deformation properties at certain locations; poor modularity, flexibility and cost-effectiveness especially due to fixing certain components or accessories directly to the rest of the cart; difficulty in interchanging components of different models or versions of carts due to shopping cart manufacturing practices of prefabricating model-specific components; limited maneuverability and leverage for a user desiring an efficient and varied displacement of the shopping cart; poor integration or removability of baskets or accessories with the rest of the cart; bulkiness or encumbering aspect of some subcomponents or accessories; and general operational deficiencies.\nThere is a need in the field for technology that can overcome at least some of the disadvantages of what is already known in the field."}
{"text": "1. (Technical Field)\nThis invention relates to improved pleated cloth and method for producing the same by providing pleating work on the cloth made of thermoplastic fiber or woven or non-woven cloth including thermoplastic synthetic fiber.\n2. (Prior Art)\nIn the prior art, cloth sewed on the front chest of a blouse for ornament is mechanically worked so that it presents a circular end forming an arc and or wavy line, after pleating work is finished for the decoration purpose.\nFor working the pleated cloth so that the hem defines an arc form and presents wavy line, it is required to erase the pleated lines at the end edge of the cloth. In the prior art method, the end of the pleated cloth are hemmed with thread (11) by utilizing a sewing machine. The cloth made according to the prior art methods is shown in FIGS. 9 and 10.\nHowever, the method as above needs another device such as a sewing machine for sewing down the end edge, besides the pleating device. The prior art method is inconvenient since the provision of the sewing machine increase the number of working processes and requires much thread indispensable for sewing down the end edge which adversely affects the appearance of the products."}
{"text": "One object of the invention is, in accordance with the preamble of claim 1, a coiling device for spring coiling machines, having exchangeable elements that contain adjusting screws, whereby each coiling pin is mounted in a receptacle so as to rotate therewith and so as to be able to rotate therein about its longitudinal axis after a clamping screw is loosened, and to be capable of moving in its longitudinal direction by means of a set screw, and whereby each coiling pin receptacle is mounted in or on a coiling pin holder, by means of which the coiling pin held in the receptacle can be pivoted in a controlled manner by means of set screws, and can be changed in position by means of an adjusting slide that mounts the coiling pin holder, by means of at least one set screw. One or two such coiling devices are preset in one finger or two finger spring coiling machines to produce helical springs from endless wire.\nAnother object of the invention is, in accordance with the preamble of claim 10, an adjusting device for coiling tools on spring coiling machines having a coiling device, for example of the above-described type, having an adjustable tool receptacle in which a coiling tool is mounted so as to be adjustable by means of a set screw in the direction of its longitudinal axis that crosses the longitudinal axis of the spring to be produced; having a tool holder that carries the tool receptacle, which is mounted thereon so as to be capable of being rotated by means of a pair of set screws supported on both sides of the tool on the tool holder; and having an adjustment slide that can be mounted laterally to the spring axis on which the tool holder is arranged so as to be capable of being pivoted about its longitudinal axis lying parallel to the spring axis by means of a bolt; whereby the tool holder is supported by means of a set screw on one side of its pivot axis on a device element that cannot slide.\nWith the aid of a suitable quick-clamp coiling pin holder and an adjusting device adapted therefor, the primary adjusting procedures should be able to be performed outside of the machine. This achieves an improved degree of utilization of the production machine, the reproducibility of a setting once attained, and machine servicing by less qualified personnel.\nThis is not possible with the presently known coiling systems, because for the most part they operate with holders that are permanently installed in the machine, and in which round coiling pins are inserted into a bore. By rotating about their own axes, these coiling pins can be brought into the necessary operating position and clamped in place. No type of adjusting aids are provided for this purpose, although it is precisely the locating of the correct angular position that has the greatest significance in the overall setting, and this requires well trained personnel.\nSubsequent fine adjustments are almost a matter of luck in this system, because the clamp must first be released for any such attempt without the previous position being retained and therefore relocatable. A reproducibility of earlier settings is for this reason not possible, although particularly with spring coiling machines this should be extremely important for reasons of the degree of utilization, because spring production cannot begin until several test springs have completed all subsequent operational processes such as start-up, setting, grinding, deburring and force testing. This often means hours of waiting time for a machine intended for production. For extremely small or extremely large spring diameters, the coiling pins, the axes of which are normally parallel to the adjusting device of the slider supporting them, must be able to pivot. There are corresponding points of rotation in the known systems, but only for an uncontrolled adjustment.\nThe object of the invention, then, is to realize exchangeable coiling pin holders that can be preset outside of the machine in which the coiling pins can move in a controlled manner in rotation about their axes, and in moving in the direction of their longitudinal axes, as well as perpendicularly thereto. In addition, it must be possible to pivot the longitudinal axes of the coiling pins, which normally run parallel to the adjustment direction of the adjustment slide carrying the coiling pin holders, in the negative or positive angular range.\nAn adjustment device must make it possible to undertake the above-referenced four adjustments outside of the machine on the basis of the values obtained from the coil ratio and the wire thickness of the springs to be produced, whereby no unintentional changes may arise through the installation in the machine.\nBecause spring steel wire can fluctuate in its hardness and therefor in its coiling behavior, a controlled subsequent adjustment must be possible in the machine."}
{"text": "Biological macromolecules such as proteins and polynucleotides have become of increasing commercial interest in medicine as pharmaceutical products. Productivity of synthetic processes is frequently limited by purification methods available. Products of biosyntheses are frequently contaminated by structurally similar impurities that must be removed before the product can be used. Chromatographic methods are typically the most effective purification methods, but the physical and chemical similarities between the desired product and the impurities frequently require laborious multiple separations.\nElution chromatography is the mode almost exclusively known and used. However, a chromatographic system may also be operated in a displacement mode, and operation in this mode can have important advantages for purification of bioproducts, particularly on a preparative and/or an industrial scale. Displacement chromatography is distinguishable from elution chromatography both in theory and in practice. In elution chromatography, a solution of the sample to be purified is applied to a stationary phase, commonly in a column. As the mobile phase is passed over the stationary phase, equilibrium is established between the mobile phase and the stationary phase. Depending on its affinity for the stationary phase, the sample species pass along the column at speeds which reflect their affinity relative to other components that may occur in the original sample.\nA modification and extension of isocratic elution chromatography is found in step gradient chromatography wherein a series of eluants of varying compositions are passed over the stationary phase.\nDisplacement chromatography is fundamentally different from elution chromatography (e.g., linear gradient, isocratic or step gradient chromatography). The displacer, having an affinity higher than any of the feed components, competes effectively for adsorption sites on the stationary phase. An important distinction between displacement and desorption is that the displacer front always remains behind the adjacent feed zones in the displacement train, while desorbents (e.g., salt, organic modifiers) move through the feed zones. The implications of this are quite significant in that displacement chromatography can potentially concentrate and purify components from mixtures having low separation factors. In the case of desorption chromatography, however, relatively large separation factors are generally required to give satisfactory resolution.\nIn displacement chromatography the eluant (i.e., the displacer) has a higher affinity for the stationary phase than do any of the components in the feed. This is in contrast to elution chromatography, where the eluant usually has a lower affinity. The essential operational feature which distinguishes displacement from elution or desorption chromatography is the use of a displacer molecule. In displacement chromatography, the column is first equilibriated with a carrier solvent under conditions in which the components to be separated all have relatively high binding. The feed solution is then introduced into the column following which the displacer is passed through the column. If the displacer and the mobile phase are appropriately chosen, the products exit the column as adjacent square waves zones of highly concentrated pure material in the order of increasing affinity of adsorption. Following the zones of purified components, the displacer emerges from the column. Finally, after the breakthrough of the displacer, the column is regenerated by desorbing the displacer from the stationary phase to allow the next cycle of operation.\nDisplacement chromatography has some particularly advantageous characteristics for process scale chromatography of biological macromolecules such as proteins. Displacement chromatography can achieve product separation and concentration in a single step unlike elution chromatography which results in product dilution during separation. Since displacement operates in the non-linear region of the equilibrium isotherm, high column loadings are possible. This allows better column utilization than elution chromatography. Finally, displacement can concentrate and purify components from mixtures having low separation factors unlike the relatively large separation factors which are required for satisfactory resolution in desorption chromatography. Displacement is thus a powerful preparative technique that can offer high production rates, resolving power and elevated yields and purity of a desired byproduct.\nThe main disadvantage of displacement chromatography, and what has limited its application in bioseparations, is the need to identify a displacer molecule for use in each separation. An effective displacer has greater affinity for the stationary phase than the bioproduct to be purified. Additionally, it should cause separation of the bioproduct from impurities on the column. Finally, it should be readily separable from the bioproduct, so that it does not become an impurity itself.\nIdentification of an effective displacer has been a laborious and tedious task. Displacer candidates are typically screened individually in column experiments using trial and error. While column experiments indicate the exact behavior of displacer molecules in the column, the time required for screening a large number of molecules is a major limitation. A technique for the high throughput screening of potential displacers would enable rapid screening of molecules generated, for example, from a combinatorial library. Screening of a large number of molecules would also provide sufficient data for a predictive QSAR model to actually direct the design of a displacer molecule for a particular bioproduct or a particular stationary phase. This would enable the identification of important properties for a particular interaction or for similar interactions on different stationary phases. Therefore, a need exists for a rapid method for screening a large number of displacer candidates."}
{"text": "1. Field\nThe present disclosure relates generally to an improved data processing system and, in particular, to a method and apparatus for allocating resources. Still more particularly, the present disclosure relates to a method and apparatus for allocating resources in network data processing systems.\n2. Description of the Related Art\nOrganizations commonly use network data processing systems in manufacturing products, performing services, internal activities, and other suitable operations. Some organizations use network data processing systems in which the hardware and software are owned and maintained by the organization. These types of network data processing systems may take the form of local area networks, wide area networks, and other suitable forms. These types of networks place the burden of maintaining and managing the resources on the organization. In some cases, an organization may outsource the maintenance of a network data processing system.\nOther organizations may use network data processing systems in which the hardware and software may be located and maintained by a third party. With this type of organization, the organization uses computer systems to access the network data processing system. With this type of architecture, the organization has less hardware to use and maintain.\nThis type of network data processing system also may be referred to as a cloud. With a cloud environment, the cloud is often accessed through the internet in which the organization uses computers or a simple network data processing system to access these resources. Further, with a cloud, the amount of resources provided to an organization may change dynamically. For example, as an organization needs more resources, the organization may request those resources.\nAs a result, organizations that use clouds do not own the hardware and software. Further, these organizations avoid capital expenditures and costs for maintenance of the computing resources. The organizations pay for the computer resources used. The organizations may be paid based on the resources actually used, such as actual processing time and storage space, or other use of resources. The organizations also may pay for fixed amounts of resources periodically. For example, an organization may pay for a selected amount of storage and processing power on a monthly basis. This usage is similar to resources, such as electricity or gas."}
{"text": "The present invention relates generally to gas turbine engines, and more particularly to a mechanical apparatus for controlling the clearance between a row of rotor blade tips and a surrounding shroud in a gas turbine engine.\nA gas turbine engine includes a rotary compressor to compress the air flow entering the engine, a combustor in which a mixture of fuel and the compressed air is burned to generate a propulsive gas flow, and a turbine which is rotated by the propulsive gas flow and which is connected by a shaft to drive the compressor. The efficiency of a gas turbine engine, such as an aircraft jet engine, depends in part on the clearance or gap between the rotor blade tips and the surrounding engine casing shroud, such as the clearance between the engine's turbine blades and the engine's turbine casing and the clearance between the engine's compressor blades and the engine's compressor casing. If the clearance is too large, more of the engine air flow will leak through the gap between the rotor blade tips and the surrounding shroud, decreasing the engine's efficiency. If the clearance is too small, the rotor blade tips may strike the surrounding shroud during certain engine operating conditions.\nIt is known that the clearance changes during engine acceleration or deceleration due to changing centrifugal force on the blade tips and due to relative thermal growth between the rotor and the engine casing. For instance, upon engine acceleration, the thermal expansion of the rotor typically lags that of the engine casing, and upon engine deceleration, the engine casing contracts more rapidly than does the rotor.\nControl mechanisms, usually of the mechanical or thermal type, have been proposed in the prior art to maintain a generally constant rotor-to-shroud clearance despite changing engine operating conditions. However, none are believed to represent the optimum design for controlling such clearance. Consequently, a need still remains for an improved apparatus for clearance control, one that will maintain a minimum clearance between the rotor blade tips and the engine casing shroud throughout the operating range of the engine and thereby improve engine efficiency by achieving more thrust with less fuel."}
{"text": "The present invention relates, in general, to specification modeling, and more particularly, to extracting a language model for a system from a natural language specification.\nWhen engineering complex systems, such as software systems, the system is commonly defined in a natural language functional specification. An important precursor to creating the system in accordance with the functional specification is the production of a model of the system. In order to create such a model, the functional specification must first be interpreted and summarized by domain experts before being converted into a model of the system in a given modeling specification language or meta-model language. Once the meta-model language model has been created it can be further converted into a model encapsulating the concepts and behaviors defined by the functional specification."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of monitoring, using the electron-beam image of a resist pattern, whether or not the exposure has been performed under an appropriate exposure condition at the time of forming the resist pattern on a semiconductor board at the lithography step in the manufacture of a semiconductor device. Also, it relates to a method of manufacturing the semiconductor device. In particular, it relates to a technology for controlling the exposure process so as to maintain the correct exposure condition, and a method of manufacturing the semiconductor device by using this condition.\n2. Description of the Related Art\nFIG. 27 illustrates a flow of the conventional lithography step.\nAt first, a resist, i.e., a photosensitive material, is coated on a board such as a semiconductor wafer with a predetermined thickness. Then, the downsize exposure of a mask pattern to the resist is performed, using an exposure apparatus (step 2050). After that, the downsize-exposed resist is developed (step 2051), thereby forming a resist pattern. Next, the size checking of the formed resist pattern is carried out, using a scanning electron microscope equipped with a size-measuring function (i.e., length-measuring SEM) (step 2052). The conventional processing contents performed using the length-measuring SEM are as follows, for example: After acquiring the electron-beam image of the region including a portion whose size-accuracy is to be strictly managed (step 2053), the resist-pattern size is measured (step 2054). Next, a judgement is made as to whether or not the size will satisfy a criterion (step 2055). If the size does not satisfy the criterion, the exposure quantity provided by the exposure apparatus is modified (step 2056, the correction quantity to the exposure quantity is represented by ΔE). In the case of, e.g., a positive-type resist, the exposure quantity is increased if the resist-pattern size is too large, while the exposure quantity is decreased if the resist-pattern width is too small. Also, in many cases, the increased or decreased quantity of the exposure quantity has been determined based on the experience and the intuition of a worker in charge.\nFIG. 28 illustrates the relationship between the resist pattern and a film pattern after being etched (cited from “Electron Beam Testing Handbook”, p. 255, Japan Society for the Promotion of Science, 98th research material of 132nd Committee on Application of Charged Particle Beam to Industry). If the etching conditions are identical, there exists a certain fixed relationship between the configuration of the resist pattern and that of the after-etched film pattern. Accordingly, in order to obtain the film pattern having a predetermined configuration, the resist pattern is also required to have the predetermined configuration. Also, when starting to manufacture a new product-type of semiconductor device or the like, prior to the introduction of a product wafer, there is performed “a condition-submitting work” for finding out a focal-point position and an exposure quantity that make it possible to acquire a predetermined resist-pattern configuration. This finding-out is carried out as follows: At first, a wafer illustrated in FIG. 12A is formed on which a pattern is pasted by being baked under a condition that the focal-point position and the exposure quantity are changed on each shot (i.e., onetime exposure unit) basis (the wafer of this kind is commonly referred to as “FEM (: Focus & Exposure Matrix) sample”). Next, in addition to the execution of the size measurement of the resist pattern on each shot basis, the wafer is cut off so as to inspect its cross-section configuration or the like. Incidentally, JP-A-11-288879 has disclosed a system for supporting the condition-submitting work. This work determines an exposure quantity (E0) and a focal-point position (F0) that allow variation margins to be taken more widely in the exposure quantity and the focal-point position. Then, based on these conditions, the exposure to the product wafer is carried out. There exist, however, some cases where a probability is increased that the resist pattern having the specification-satisfying configuration cannot be acquired. This is because various process variations (e.g., a change in the photosensitivity of the resist, a film-thickness variation in an anti-reflection film under the resist, and drifts in the respective types of sensors of the exposure apparatus) decrease the variation margins under the conditions (E0, F0) determined in the condition-submitting work. Detecting these process variations is a role to be played by the above-described size measurement (i.e., the step 2). In the prior arts, the attempt has been made to compensate, by the correction to the exposure quantity, for the size change in the resist-pattern configuration caused by the process variations.\nIn the prior arts, in order to detect and take a countermeasure against the process variations, the following method has been employed: The size value of a line width or the like is inspected, using the length-measuring SEM. Then, if the size value does not satisfy a criterion, the exposure quantity is corrected. This method, however, is accompanied by the following 1st to 3rd problems:\nAs the first problem, the length-measuring SEM usually observes, from directly above, the resist-pattern configuration formed on the semiconductor board. Also, the length-measuring SEM measures the size value of the edge configuration of the resist pattern that has appeared on the electron-beam image acquired in this state. There exists, however, a variation in the resist-pattern configuration which cannot be measured by conventional measuring techniques. Concretely, this variation means a variation in the resist-pattern's edge configuration caused by the variation in the focal-point position at the time of the exposure. The conventional measuring techniques find it impossible to detect this variation. The cross-section configuration of the resist pattern formed on the semiconductor board is a substantially trapezoidal configuration. Concerning the signal intensity of secondary electrons detected by the scanning electron microscope as a signal generated from the sample, the signal intensity from the inclined portion is stronger than the one from the flat portion. As a result, as illustrated in FIG. 29A, the signal waveform becomes a one that has the peaks at the locations corresponding to the edges of the trapezoid. In the size measurement by the length-measuring SEM, as illustrated in, e.g., FIG. 29B, the following method or the like is employed: Straight lines are applied to the outer-side portions and the base portions of the peaks so as to determine each intersection point of the respective 2 straight lines, then defining, as the line width, the distance between the right and left intersection points. FIG. 30 is a graph for indicating how the line width will change if the exposure quantity and the focal-point position are changed, where the focal-point position is set up in the transverse-axis and the line width is plotted on each exposure-quantity basis (i.e., e 0 to e 8). The exposure quantity lies in the range of e 0<e 1< . . . <e 8. Moreover, there exists a relationship that the line width becomes smaller as the exposure quantity grows larger (, which is in the case of a positive-type resist. This relationship becomes opposite in the case of a negative-type).\nConsequently, inspecting the line width makes it possible to detect the variation in the exposure quantity. As is apparent from the same graph, however, the line width does not change greatly with respect to the change in the focal-point position. In particular, in proximity to the correct exposure quantity (i.e., e 4), the line width scarcely changes if the focal-point position is changed. Accordingly, the inspection of the line width does not make it possible to detect the variation in the focal-point position. Meanwhile, even if the line width does not change, the cross-section configuration of the resist pattern changes if the focal-point position is changed as is illustrated in FIG. 30B. As described earlier, the change in the cross-section configuration of the resist pattern exerts the influence on the film pattern after being etched. As a consequence, the use of the conventional measuring techniques, which are incapable of detecting the variation in the focal-point position, may result in even a danger of producing a configuration failure of the after-etched film pattern in large quantities.\nAs the second problem, there exists a problem that the correction of the exposure quantity alone, naturally, finds it impossible to deal with the case where the focal-point position has been deviated. In the case of, e.g., the situation A in FIG. 30A, since the line width is larger than the normal line width, the processing of increasing the exposure quantity will be carried out based on the measurement result of the line width. This processing, however, makes no correction to the deviation in the focal-point position, thereby simply bringing about the situation B in FIG. 30A and never implementing the normal cross-section configuration of the resist pattern. This, accordingly, may again result in even the danger of producing the configuration failure of the after-etched film pattern in large quantities.\nAs the third problem, in the above-described prior arts, there also exists a problem of being unable to acquire information for indicating the process variations in a quantitative manner. Here, this information is needed in order to maintain the normal exposure process. In accompaniment with the microminiaturization of a semiconductor pattern in recent years, the variation allowable range of the exposure quantity and that of the focal-point position have become exceedingly small. For example, it is requested that, when the pattern rule is smaller than 180 nm, the variation range of the exposure quantity is controlled to become smaller than ±10% of the size value, and that of the focal-point position is controlled to become smaller than ±0.2 to 0.3 μm. The implementation of these variation allowable ranges requires the information for indicating the process variations in a quantitative manner, i.e., an accurate quantification of the variation quantities. This accurate quantification is such that the deviation in the exposure quantity is equal to such-and-such mJ, and the deviation in the focal-point position is equal to such-and-such μm. In the prior arts, the detection of the deviation in the focal-point position is not performed at all. Moreover, it cannot help saying that the detection of the deviation in the exposure quantity is inaccurate. The reason for this is as follows: Despite the fact that, in general, the variation in the focal-point position also changes the line width, the size variation quantity caused by the variation in the focal-point position has also been compensated for by the adjustment of the exposure quantity. Consequently, the use of the prior arts finds it impossible to maintain the normal exposure process."}
{"text": "A current method for the large scale collection and concentration of solar energy is the use of heliostat mirrors to reflect the sun's rays to a central receiver. By utilizing multiple heliostat mirrors, each one reflecting to a common point, concentration of solar energy is achieved.\nIn known systems, heliostat mirrors may be in a fixed position surrounding the tower. The mirror surfaces are typically controlled in two degrees of motion to position the surface of the mirror with respect to the tower. Each heliostat mirror has a control system which tracks the motion of the sun with respect to the centrally located receiver. The mirror is on a fixed base and the reflective surface of the mirror is continuously moved to maintain the solar reflection from the surface of the mirror onto the receiver. The purpose of positioning the heliostat mirrors being to reflect and direct the sun's rays to a designated central collection point, known as a central target receiver or a power tower. In order to accomplish this, the heliostat mirror requires a surface area of reflective mirror, two axes of motion, a servo motor for each axis of motion and a control system for positional calculation and motion control of the two axes.\nIn some installations, the array of heliostats encompasses a full 360 degree area around a receiver. This is a very expensive installation due to the large number of heliostats required in order to populate the array and surround the receiver. Also, because the position of the sun is changing throughout the day, at any point in time more than half of the array has a very low performance due to their fixed placement within the array.\nOther installations have attempted to address the issue of heliostat cost by reducing the number of mirrors. A reduced number of mirrors are controlled as a gang in that the heliostat array is a single movable unit that is positioned on a track. The array then rotates concentrically about the receiver on a horizontal plane of motion. The mirrors in the array are not only controlled as a single unit about the receiver, but are also movable as a single unit about a vertical axis to rotate their surface position with respect to the receiver. These installations are not practical for an installation having a large radius from the central tower. Also, large arc segments of mirrors will not target the receiver with adequate accuracy.\nFurthermore, known concentric designs do not always take into account the requirement of a full 360° for rotation about the central receiver. In global implementations, particularly for solar tracking in latitudes between the Tropic of Cancer and the Tropic of Capricorn which comprises a large global market segment for solar thermal applications, known concentric designs may not capable of tracking the sun. For this special case, the solar path is not always southerly or northerly, but both depending on the day of the year. In fact, the solar path may be north or south of the central receiver and at one point, will be directly overhead. Without the capability of a full circle of motion for the array, known concentric designs will not operate within this equatorial band.\nAnother major factor in a solar thermal concentration system is the environmental impact of the installation. In known systems, a large area of land is necessary and that land must be dedicated to the sole purpose of solar thermal power collection. Additionally, the typical ground preparation for installation involves invasive grading and leveling practice which completely destroys the indigenous habitats and ecology of the terrain. Furthermore, the dedicated area is permanently covered by the mirror array, limiting ground exposure to natural heal and light provided by the sunlight, which negatively impacts the environment.\nThere is a need for a solar thermal concentration system and method that is cost effective to install and maintain, fully utilizes the collection possibilities as the earth orbits the sun, and has minimal adverse impact on the environment."}
{"text": "1. Field of the Invention\nThe present invention relates to a multi-wavelength laser diode, more particularly, which can improve a reflective layer structure in order to achieve a high reflectivity in the entire visible light range.\n2. Description of the Related Art\nLaser diodes are used as an optical pickup for various storages in order to produce light in response to application of electric current.\nA laser diode needs a high reflectivity mirror or reflective face on one end of an oscillating layer in order to realize high power laser beams, in which producing such a reflective face requires coating a high reflectivity thin film. Various studies have been made so far in order to produce such a high reflective thin film.\nFIG. 1 is a perspective view of a conventional laser diode, FIG. 2 is a cross-sectional view taken along the line II-II in FIG. 1, and FIG. 3 is another cross-sectional view taken along the line III-III in FIG. 1.\nReferring to FIGS. 1 to 3, a laser diode 10 includes an oscillating structure 12 for generating light, p- and n-metal layers 26 and 28 coated on the top and underside of the oscillating structure 12, respectively, a reflective layer 30 of a high reflectivity formed on one end of the oscillating structure 12 and an anti-reflective layer 40 formed on the other end of the oscillating structure 12 opposite to the reflective layer 30.\nThe oscillating structure 12 includes a lower cladding layer 16, an active layer 18, an upper cladding layer 20 and an upper capping layer 24 formed in their order on a semiconductor substrate 14. The upper cladding layer 20 has a ridge 20a formed in a central upper portion thereof so that inactive regions 22 such as current blocking layers are formed at both sides of the ridge 22a. \nWhen a forward voltage is applied to the p- and n-metal layers 26 and 28 of the laser diode 10 having the above oscillating structure 12, electrons and holes are introduced along vertical routes into the active layer 18 where the electrons and holes are combined to emit a laser beam, i.e., photons corresponding to the energy band gap of the active layer 18.\nThe reflective layer 30 is formed on one face of the oscillating structure 12 adjacent to one end of the ridge 22a and the anti-reflective layer 40 is formed on another face of the oscillating structure opposite to the reflective layer 30 so that a cavity resonator formed between the reflective and anti-reflective layers 30 and 40 can generate laser oscillation.\nThe reflective layer 30 is made of dielectric material, and includes at least 7 pairs of Al2O3 and Si3N4 layers which are deposited one atop another to raise reflectivity.\nHowever, a process for forming the reflective layer 30 of multilayer dielectric material spends a long time period while producing poor reliability. Although the multilayer dielectric layer can achieve excellent reflectivity with respect to a single wavelength, there is a critical drawback in that it cannot achieve a uniform reflectivity with respect to a multi-wavelength band.\nThat is, the thickness of each dielectric layer is determined by λ/4n, in which λ is the wavelength of a laser beam to be reflected, which is emitted from an oscillating layer including the active layer 18 and the upper and lower cladding layers 16 and 20, and n is the reflectivity of corresponding dielectric material.\nReflectivity examples of conventional dielectric reflective layers with respect to wavelengths are shown in FIGS. 4 and 5, in which FIG. 4 illustrates a reflectivity curve of a dielectric reflective layer with respect to wavelengths where a corresponding laser beam has a wavelength of 715 nm, and FIG. 5 illustrates a reflectivity curve of a dielectric reflective layer with respect to wavelengths where a corresponding laser beam has a wavelength of 780 nm.\nThe graph shown in FIG. 4 has the highest reflectivity of about 92% at 715 nm. However, it is seen that the reflectivity significantly drops at a wavelength range of 650 nm or less. The graph shown in FIG. 5 also has the highest reflectivity of about 92%. It can be observed, however, that the reflectivity also significantly drops at a wavelength range of 700 nm or less.\nAccordingly, such a reflective layer based upon dielectric material can have a high reflectivity only in a specific wavelength range, but cannot maintain the high reflectivity in a wide wavelength range."}
{"text": "This invention relates to a medical device, and more particularly, to a highly safe medical device. More illustratively, this invention is directed to a highly safe blood bag system having an excellent blood- or blood component-storability.\nSoft bags fabricated from polyvinyl chloride are recently used to store blood for transfusion instead of rigid glass containers mainly in view of the reduced damage of erythrocytes during the storage. Use of such polyvinyl chloride soft bags has increased also for their superior workability, flexibility, transparency, resistance to vapor permeation, heat resistance, and the like. Conventional polyvinyl chloride resins which has been used for fabricating such blood storage bags contain from about 30 to 60 parts by weight of a phthalate for plasticizing purpose. Such phthalate elutes into blood during the blood storage, and it has been found out that the phthalate in blood protects membrane of various cells of the blood, in particular, erythrocytes. Despite such positive effects, the phthalate eluted into the blood may induce safety problems upon its introduction into the body by blood transfusion."}
{"text": "Highly fluorinated ion-exchange polymers, such as the sulfonyl type disclosed in U.S. Pat. No. 3,282,875, are used in membrane form as separators in electrochemical cells. The polymers are also useful as acid catalysts. These applications first used melt-fabricated, i.e. melt-processed, shapes such as films and pellets. Because the ion-exchange polymers are difficult to melt-process in the ionic form, fabrication is carried out on polymer in a melt-processible precursor form, and the fabricated article is then hydrolyzed to convert the polymer to the ionic (also referred to as the ion-exchange) form. Later, methods were discovered for making liquid compositions of highly fluorinated ion-exchange polymers: U.S. Pat. Nos. 4,433,082 and 6,150,426. From such liquid compositions, ion-exchange membranes can be made by film-casting techniques. Catalyst can be made by coating liquid compositions on inert substrates. Liquid compositions have also found use in making electrodes for fuel cells.\nDissolution of the above polymers in their ionic forms requires solvents and temperatures sufficient to overcome the forces that hold the polymer together in the solid state. These forces include the polar attractions of the ionic groups for one another. Polar solvents such as water and alcohol can solvate the ionic groups of the polymer, weakening their interaction, and promoting dissolution. Heat further weakens intermolecular attractions.\nLiquid compositions are typically made by putting highly fluorinated ion-exchange polymer that is in the sulfonic acid form, in alcohol, in water, or in aqueous alcohol, and heating the combination to achieve the dissolution of the polymer. The temperatures necessary are generally 220° C. or higher. These temperatures are above the boiling point of the solvent at atmospheric pressure and therefore the dissolution is conducted in an autoclave. Higher temperatures are necessary with water. Alcohol is a better solvent, and lower temperatures can be used. However, the formation of side products such as ether and olefin through the reaction of the alcohol with the strongly acid polymer contributes to the development of pressure during dissolution. Water-alcohol, i.e. aqueous alcohol, is effective at lower temperatures than are necessary with water alone, and also does not develop pressures so high as occur with alcohol alone. Nevertheless, though reduced in quantity, side products still form, adding to reaction pressure, requiring separation from the liquid composition and disposal, and resulting in loss of solvent. New solvents are needed that are effective at lower temperatures and that produce less side product."}
{"text": "There is a particular need for preventing the accidental discharge of a weapon, such as a hand weapon, by an unauthorized person.\nThere are many problems that can occur with the use of weapons, automatic weapons and handguns. Guns are usually provided with a \"safety,\" which is intended to prevent an accidental discharge by inhibiting the firing of the gun when the safety button or lever is switched on. The use of a \"safety lever\" does not help in preventing an unauthorized individual from using another person's gun.\nOne example of this type of problem is when a child finds a hidden gun in his home, where its unintentional use may cause serious injury or even the death of another child. Still another example of an accidental misuse of a weapon is where a criminal wrests a police officer's gun from him and uses it on him.\nMeasures to prevent the unauthorized use of a handgun include many procedures or devices, such as locking up the gun in a safe place, or by using devices like trigger locks or barrel plug locks to render the gun inoperative. While these measures prevent the unauthorized use of the gun, they also hinder the use of the gun by its rightful owner since it takes time to unlock and prepare the gun for use. These mechanisms would be useless to police officers, especially in a situation that was previously described.\nIt is an object of this invention to provide a new method and safety system that will prevent a gun from firing unless it is being held by the gun's owner. Any other unauthorized person who has not registered with the safety device will not be able to fire the weapon.\nIt is another object to provide a transponder means external to the weapon, such as a finger ring, as a means for exchanging data between the intended weapon and the registered user of the weapon having the finger worn ring.\nYet, another object of this invention is to provide a method of registration of a handgun with its accompanying finger ring that provides a verifiable identification code."}
{"text": "1. Technical Field\nThe present invention relates to composite golf ball-hitting mechanisms, and more particularly, to a composite golf ball-hitting mechanism for providing golf tees which come in different sizes available to golfers for selection.\n2. Description of Related Art\nOne of the popular sports nowadays is golf. In addition to golf clubs with various functions, golfers need golf tees each for holding a golf ball for the first stroke of each hole. Golf tees of different heights are required to propel the golf balls different distances. Referring to FIG. 1 and FIG. 2, a conventional golf tee 1 essentially comprises a shaft portion 11 with two ends, an apex 12 disposed at one of the two ends, and a dish 13 disposed at the other one of the two ends. After the apex 12 has been inserted into the lawn of a golf course, a golf ball 2 is put on the dish 13 and then hit. However, the size of the golf tee 1 is the key to the reach of the golf ball 2. For this reason, not only must a golfer carry the golf tees 1 of different sizes simultaneously, but the golfer must also insert the golf tee 1 into the lawn of the golf course before hitting the golf ball; as a result, not only does the golf tee 1 lack ease of use, but the golfer is likely to hit and fling the golf tee 1 when hitting the golf ball, and thus eventually the golfer has to gather the golf tee 1 located afield. Last but not least, the golf tee 1 has the apex 12; hence, not only does the flinging golf tee 1 pose a threat to sports safety, but the storage of the golf tee 1 is also difficult."}
{"text": "Wireless communications devices such as the BlackBerry® by Research in Motion Limited provide a variety of useful functions, such as voice communication, e-mail and Web browsing. Of growing popularity are mapping applications that take advantage of the device's wireless link and Web capabilities to download and render maps on the device. When coupled with a GPS receiver, either embedded as a GPS chipset or externally (e.g. Bluetooth™ connected), these mapping applications can be used for (or modified for use with) navigation and other location-based services (LBS).\nNavigation technology, be it on dedicated GPS navigation units or on multi-function wireless handheld devices, has evolved rapidly in recent years from providing a graphical representation of a route superimposed on a map and/or a list of driving directions to full-blown turn-by-turn navigation where real-time instructions are presented to the user visually and/or audibly (using text-to-speech technology). In addition to current position and route information, current mapping and navigation software often provide searchable databases of points of interest (POI) such as landmarks, hotels, restaurants, cafes, gas stations, etc., which the user can optionally have displayed on the map. These POI can be either stored in onboard memory or retrieved wirelessly on-demand from a map server. Although street names and points of interest help a user navigate, it would be highly desirable to provide more detail about the device's current position to further facilitate navigation."}
{"text": "In recent years, Wavelength Division Multiplexed (WDM) and Dense Wavelength Division Multiplexed (DWDM) optical transmission systems have been increasingly deployed in optical networks. Although DWDM optical transmission systems have increased the speed and capacity of optical networks, the performance of such systems, especially those providing high bit rates, has traditionally been limited by various factors such as chromatic dispersion and the non-linearity in an optical fiber's refractive index, which can cause spectral broadening of optical pulses and degrade the transmission of high speed optical signals. Because such optical signal degradation tends to accumulate along transmission paths, chromatic dispersion and non-linearity can significantly limit the transmission distance of high speed optical signals.\nChromatic dispersion refers to the fact that different wavelengths of light pass through an optical fiber at different speeds, thereby causing a pulse of light propagating through the optical fiber to broaden. Chromatic dispersion is an inherent property of the glass fiber and arises from two sources, the material properties of the glass and the waveguide structure (i.e., the refractive index profile of the glass fiber). For example, a single mode fiber through which light propagates at a wavelength of 1550 nm operates in the so-called anomalous dispersion regime. In the anomalous dispersion regime longer wavelengths (i.e. lower frequencies) of light travel at slower speeds than shorter wavelengths (i.e. higher frequencies) of light. As a result, the different wavelengths in pulse of light will be broadened as it travels through the single mode optical fiber.\nSeveral solutions have been proposed to mitigate the effects of dispersion in transmission fibers. One technique involves the use of a compensating optical fiber having an appropriate length and which has a dispersion that is opposite to the dispersion characteristic of the transmission fiber. As a result, the dispersion in the transmission fiber is substantially canceled by the total dispersion in the compensating fiber. The dispersion of the compensating fiber is generally selected to be much greater in magnitude than the dispersion of the transmission fiber. In this way the length of the dispersion compensating fiber may be much less than the length of the transmission fiber. For example, to compensate for dispersion in a single mode fiber 100 km in length, a typical dispersion compensating fiber may need to be about 10 km in length. Such fiber is generally wound on a mandrel and provided for use as a dispersion compensating module.\nSince high dispersion fiber has a relatively small core area (e.g., about 25 μm2 for a −100 ps/nm-km fiber), the optical power density in the dispersion compensating fiber will be relatively high, which leads to signal degradation arising from non-linear interactions. Accordingly, there is limit to the optical power that can be directed into the dispersion compensating fiber if nonlinear penalties are to be avoided. For example, it has been determined that a safe power level for light launched into one common dispersion compensating fiber is about 0 dBm/wavelength. Thus, for a WDM system employing 40 channels or wavelengths, the maximum permissible power that can be launched into the dispersion compensating fiber while avoiding non-linear penalties is about 16 dBm.\nIn many cases WDM networks are often initially deployed at less than their maximum capacity. That is, a system designed to transmit 40 channels or more, for instance, initially may be more lightly loaded with only 2, 4, or 8 channels. As demand increases, network capacity can be increased by increasing the number of wavelengths that are used. In addition, the number of wavelengths that are used may be changed dynamically as demand increases or decreases. Such dynamic networks will need to dynamically control the input power to the dispersion compensating fiber to avoid nonlinearities.\nAccordingly, it would be desirable to provide a method and apparatus for automatically adjusting the optical power directed to a dispersion compensating fiber so that nonlinear interactions are avoided."}
{"text": "1. Field of the Invention\nAn embodiment relates to a light emitting apparatus and a light unit.\n2. Description of the Related Art\nIn recent years, a light emitting apparatus and a lighting system using a light emitting device as a light source have been widely used. The light emitting diode (LED) may configure the light source that generate light by using compound semiconductor materials such as GaAs-series, AlGaAs-series, GaN-series, InGaN-series, and InGaAlP-series.\nMeanwhile, since the emission characteristics of the LED are deteriorated due to heat, the light emitting apparatus and the lighting system that have the LED need to be designed to rapidly emit heat emitted from the LED to the outside."}
{"text": "Water purification apparatus for use in laboratories and healthcare facilities are well known. Generally they involve the reduction and/or removal of contaminants and impurities to low levels. They typically contain a variety of technologies that remove particles, colloids, ionic species and organic substances and/or molecules.\nWater purification apparatus generally allow manual variation by a user of the volume of water to be dispensed, such as from drops to liters, and/or the speed (i.e. flow rate) of water dispense, again from drop-wise to a number of liters per minute.\nHowever, as well as manual dispensing based on a user physically operating a tap and manually gauging the volume of water being dispensed, frequently water purification apparatus provide for ‘automatic dispensing’, such as on a time basis or measured basis; optionally following a pre-programmed or pre-stored dispense profile as discussed in our WO2007/015048A1.\nIn order to accommodate such a range of dispensing options, the control of the water dispense outlet of water purification apparatus is often ‘electronic’, to automate this activity. Typically, a simple on/off solenoid-controlled valve is used. However, solenoids are limited in that the internal flow rate of the water dispense apparatus, and hence the flow rate of the dispense of the purified water stream, is conventionally set at a single and high value, and this does not allow for the easy dispense of slow or small volumes, particularly into smaller water collecting vessels such as beakers and flasks.\nU.S. Pat. No. 5,925,240 describes a water treatment system having dosing control for providing treated water that can be accurately dispensed. In one embodiment, its controller converts an input signal to a motor control signal that causes the motor to operate at a speed which drives the pump and causes the flow of water at the outlet valve to correspond to the desired flow rate. In another embodiment of U.S. Pat. No. 5,925,240, the controller controls the outflow cross-section of a valve and causes the flow of water at the valve to correspond to the desired flow rate. However, both such embodiments involve increasing costs for additional solenoids and additional hydraulic complexity.\nIn order to overcome such limitations, but retain a good degree of dispense control, voltage or frequency dependent solenoids can be used, whereby adjustment of the coil, voltage or frequency adjusts the position of the armature, enabling some flow adjustment through such solenoid-controlled valves. However, such solenoids are still limited at very low and/or very accurate (such as drop wise) flows that can be desired or required when very precise volumes are to be dispensed. This is because on initial lifting of the valve diaphragm to start a dispense, the change in pressure on the valve diaphragm affects the delicate electronically-floating armature, leading to its incorrect positioning and therefore incorrect flow control.\nA further disadvantage of solenoids is that in energising the activation coil of a solenoid over a long period, (such as during drop wise dispense for very accurate volume dispense), the heated coil causes undesirable localised heating of the passing purified water stream to occur, affecting the properties of the purified water stream."}
{"text": "Concerns about the environmental impact of fluorocarbons and hydrofluorocarbons on the environment call for expanded uses of environmentally sound refrigerants, such as water or ammonia. For commercial chiller applications in the capacity range of 15-10,000 refrigeration tons with little heating requirements, water is usually the preferred refrigerant due to its non-flammable and benign nature.\nAqueous absorption fluid cycles taking advantage of such refrigerants have been known and used for many decades. Single effect and various two-stage designs are commercially used in many countries. However, the increasing concern about the overall CO.sub.2 generation in the process of converting fossil fuels into energy used for air conditioning or refrigeration calls for higher energy conversion efficiencies than currently obtained with single stage (COP=0.6 to 0.8) or two-stage (COP=0.9 to 1.25) absorption equipment.\nU.S. Pat. No. 4,732,008 teaches the use of two single stage cycles coupled to achieve three refrigeration effects. The independent loops carry different absorption fluids with the lower stage using fluids such as aqueous LiBr solutions. However, the needed upper stage fluid crystallization and vapor pressure suppression properties require the use of a different fluid. Performance estimates lead to COPs of 1.5 to 1.7 if suitable upper stage fluids are used."}
{"text": "1. Field of the Invention\nThe present invention relates to a seat belt system which is provided with a take-up apparatus for taking up a webbing to be placed around an occupant of an automobile in the form of layers around a take-up shaft.\n2. Description of Prior Art\nA seat belt system mounted in an automobile is provided with a take-up apparatus which takes up one end of the webbing that is used to secure an occupant and stores the webbing in the form of layers around a take-up shaft so that the occupant can pull the webbing out of the take-up apparatus and place the webbing.\nSuch a take-up apparatus is provided with an inertia lock mechanism which stops the webbing from being pulled out during emergency by inhibiting the rotation of the take-up shaft in the direction of pulling out of the webbing. Thus, the occupant is secured by the webbing at the time of a vehicle emergency.\nHowever, when the rotation of the take-up shaft in the direction of pulling out of the webbing is inhibited during an emergency in this way, a high degree of tension is then applied to the webbing due to the inertia force of the occupant, and thus tightening occurs in the roll of webbing wound around the take-up shaft in layers. This causes the webbing to be pulled out from the take-up shaft even after the rotation of the take-up shaft in the direction of pulling out of the webbing has been inhibited."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of and a device for acoustooptically forming images by applying light to a recording medium having different optical densities representative of an image to generate an image-dependent pressure wave.\n2. Description of the Prior Art\nVarious processes are used in the art of forming or reproducing images recorded on recording mediums. Cameras for photographing radiation images are also known. One conventional radiographic camera is an X-ray camera using an X-ray photographic film as the recording medium. However, the formation of an image based on the optoacoustic effect is not known in the art."}
{"text": "The present invention relates to the synthesis of novel morpholine derivatives and the evaluation of their hypocholesterolemic, hypolipidemic and antioxidant activity. Especially, the present invention relates to the synthesis and pharmacochemical evaluation of morpholine derivatives of the formula I \nin which R1=CH2CH3, R2=CH3, R3, R4=H, R5=C6H5 (compound 1) or R1=CH2CH2CH2ONO2, R2=CH3, R3,R4=H, R5=C6H5 (compound 2) or R1=H, R213 R3=(CH2)4, R4=H, R5=C6H5 (compound 3) or R1=CH2CH2CH3, R2xe2x80x94R3=(CH2)4, R4=H, R5=C6H5 (compound 4) or R1=CH2CH2CH2ONO2 R2xe2x80x94R3=(CH2)4, R4=H, R5=C6H5 (compound 5) or R1=H, R2=CH3, R3xe2x80x94R4=(CH2)4, R5=C6H5 (compound 6) or R1=CH2CH2CH3, R2=CH3, R3xe2x80x94R4=(CH2)4, R5=C6H5 (compound 7) or R1=CH2CH2CH2ONO2, R2=CH3, R3xe2x80x94R4=(CH2)4, R5=C6H5 (compound 8) or R1=CH2CH2CH2ONO2, R2=CH3, R3, R4=H, R5=H (compound 9) or R1=H, R2=p-NO2xe2x80x94C6H4xe2x80x94CH2CH2, R3, R4=H, R5=C6H5 (compound 10).\nIt is well known that the pathogenesis of atheromatosis is related with a lot of factors, the most important of them are:\noxidative modification of the low density lipoproteins (LDL)\nincreased levels of cholesterol and LDL-cholesterol in blood\nincreased levels of triglycerides in blood\ndecreased levels of HDL-cholesterol in blood\nthrombogenesis, endothelial injury and haemodynamic factors.\nEspecially, the oxidative modification of LDL appears to be the most risk atherogenic process, which induces inflammatory and apoptotic mechanisms and finally the formation of foam cells and fatty streaks.\nAt this field of research, the most of the synthetic compounds posses hypocholesterolemic or hypolipidemic or antioxidant activity. Their are no reports about compounds able to decrease total cholesterol, triglycerides and LDL-cholesterol in blood in combination with antioxidant activity.\nConsidering the above we think that it would be interesting to design, synthesise and evaluate compounds with hypocholesterolemic, hypolipidemic and antioxidand activity.\nThe novel derivatives we are reporting are able to decrease total cholesterol levels in plasma significantly. Also, the novel derivatives are able to decrease triglycerides and LDL-cholesterol levels in plasma.\nThe novel derivatives, except their hypolipidemic activity posses also potent antioxidant activity. We think that the combination of the hypolipidemic and the antioxidant activity is necessary for the prevention and treatment of ahteromatosis more effectively.\nThe present invention relates to the synthesis and the evaluation of antioxidant, hypocholesterolemic and hypolipidemic activity of novel morpholine derivatives of the general structure (I).\nSynthesis\nThe 2-hydroxy-morpholine derivatives 3, 6 and 10 are synthesised by the reaction of the appropriate aminoalcohol (22 mmol) and the 2-bromo4-phenylacetophenone or the 2-bromoacetophenone (10 mmol) in ether and acetone for 15 hours at room temperature. After washing with saturated solution of sodium chloride, drying with potassium carbonate, evaporation in vacuo and neutralisation with hydrochloric acid 10% in ether, the 2-hydroxy derivatives are obtained and recrystalised from acetone and ether.\nThe 2-alkoxy derivatives 1, 4 and 7 are synthesised by the reaction of the respective 2-hydroxy derivative with the appropriate alcohol, in acid medium and reflux.\nCompounds 2, 5, 8 and 9 are synthesised by the reaction of the respective 2-hydroxy derivative with the 3-bromopropanol in acidic medium and reflux. The 2-(3-bromopropoxy) derivatives then react with silver nitrate in acetonitrile and reflux. \nEvaluation of the Antioxidant Activity\n[Method]\nThe evaluation of the antioxidant activity of the novel compounds was performed in vitro by peroxidation of rat hepatic microsomal membrane lipids. Lipid peroxidation was induced by the Fe2+/ascorbic acid system. Rat hepatic microsomes were heat-inactivated (90xc2x0 C. for 90 s). A fresh solution of ascorbic acid (0.2 mM) in Tris-HCl buffer (pH 7.4) and the tested compounds, dissolved in DMSO to give final concentrations of 1 mM to 0.1 mM, were added to microsomes Equal amount of buffer was added to the control samples. The reaction was initiated by addition of FeSO4 (10 xcexcM). The mixture was incubated at 37xc2x0 C. for 45 min. Aliquots were taken at various time intervals and lipid peroxidation was assessed by spectrophotometric determination of the 2-thiobarbituric acid reactive material at 535 nm.\n[Results]\nAll tested compounds inhibited lipid peroxidation 100% at 1 mM and this action is maintained even at 0.1 mM.\nEvaluation of the Hypolipidemic Activity\n[Method]\nThe hypocholesterolemic and hypolipidemic activity of the synthesised compounds was performed by the inhibition of the hyperlipidemia induced by the Triton WR 1339 in rats. Triton WR 1339 as an intraperitoneal (i.p.) injection acts due to induction of HMG-CoA reductase. Using male Fischer rats (230-280 g) in group of five (test group), we injected a solution of 200 mg/Kg Triton WR 1339 i.p. and the same time 28 to 56 xcexcmol/Kg of the tested compounds and probucol (as reference), suspended in aqueous solution was injected i.p. Control group was treated i.p. with Triton WR 1339 and 10 ml/Kg of the vehicle of the tested compounds i.p. After 24 hours, blood was taken from the celiac aorta and collected to heparinised tubes. Blood was centrifuged for 15 minutes at 3000 rpm and plasma was selected for the determination of total cholesterol (TC), LDL cholesterol (LDL-C) and triglycerides (TG) using commercial kits.\n[Results]\nOur results indicate that the tested compound decrease total cholesterol levels in plasma even by 54%, triglyceride levels even by 49% and LDL-cholesterol level even by 51% and all compounds were more potent than probucol (table 1).\nThe above results indicate that the novel compounds posses antioxidant and potent hypolipidemic and hypocholesterolemic action. This combination of antioxidant and hypolipidemic properties may be useful against atheromatosis."}
{"text": "1. Field of the Invention\nThe present invention relates to Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass and Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallizable glass therefor.\n2. Description of the Related Art\nConventionally, Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass has been used for front glass panels in kerosene heaters, wood stoves, and the like, substrates for high-technology products such as color filters and image sensors, setters for baking electronic devices, shelf boards for microwave ovens, top plates for electromagnetic cooking devices, window glass for fire doors, and the like.\nFor example, Japanese Examined Patent Publication (JP-B) No. S39-21,049, Japanese Examined Patent Publication (JP-B) No. S40-20,182, and Japanese Unexamined Patent Publication (JP-A) No. H01-308,845 disclose Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass having a xcex2-quartz solid solution represented by the formula Li2Oxe2x80x94Al2O3.nSiO2, in which nxe2x89xa72, or a xcex2-spodumene solid solution represented by the formula Li2Oxe2x80x94Al2O3.nSiO2, in which nxe2x89xa74, both of which are produced in the crystallized glass as a main crystal.\nThe Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass mentioned above has high mechanical strength and superior thermal characteristics, such as a low coefficient of thermal expansion. In addition, since the crystal produced in the Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass is alterable by changing heating conditions in a crystallization process, transparent crystallized glass (a xcex2-quartz solid solution (produced) and white and opaque crystallized glass (a xcex2-spodumene solid solution is produced) can be produced from the same composition of raw glass materials. Consequently, in accordance with the application, Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass suitable for the application can be selectively produced.\nWhen this kind of crystallized glass is produced, it is required to heat up to a high temperature above 1,400xc2x0 C. for melting raw glass materials. Accordingly, as a fining agent added to one batch for glass production, arsenic oxide (As2O3) is used which is capable of evolving a great deal of gas for fining during melting at a high temperature. In melting in a batch production system, As2O3 contained in raw glass materials is oxidized at 400 to 500xc2x0 C. so as to form As2O5, and the As2O5 is again reduced at 1,200 to 1,800xc2x0 C. to form As2O3, whereby oxygen gas is evolved. The oxygen gas thus evolved diffuses into bubbles in the glass, and the bubbles enlarge, facilitating rising thereof, whereby the bubbles are removed. Since As2O3 has the fining effect described above, As2O3 has been widely used as a fining agent for glass. In particular, As2O3 is a very effective fining agent for Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass that must be subjected to melting at a high temperature.\nIn addition to the fining effect, As2O3 has an effect of promoting crystallization. Accordingly, it is understood that As2O3 is an essential component in order to obtain desired characteristics of the crystallized glass described above.\nHowever, As2O3 is highly toxic, and hence, environmental pollution may occur during a manufacturing process for glass, disposal of waste glass, and the like. Accordingly, it is desirable to reduce the amount of As2O3 used. However, when the amount of As2O3 is simply decreased, the effects on fining and crystallization are degraded, and as a result, clarity and glass characteristics, which are equivalent to those obtained if the amount of As2O3 were not decreased, cannot be obtained.\nAccordingly, it is an object of the present invention to provide Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass having clarity and glass characteristics equivalent to or superior to those of conventional crystallized glass, even when the amount of As2O3 used is decreased, and to provide Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallizable glass as a raw glass material therefor.\nAccording to one aspect of the present invention, there is provided a Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallized glass which consists essentially of, by weight, 60 to 75% SiO2, 17 to 27% Al2O3, 3 to 5% Li2O, 0 to 0.9% MgO, 0 to 0.9% ZnO, 0.3 to 5% BaO, 0 to 3% Na2O, 0 to 3% K2O, 0 to 4% TiO2, 1 to 4% ZrO2, 0 to 4% P2O5, 0.05 to 2% Sb2O3, 0 to 0.9% MgO+ZnO, and 0 to 4% Na2O+K2O.\nAccording to another aspect of the present invention, there is provided a Li2Oxe2x80x94Al2O3xe2x80x94SiO2 crystallizable glass which consists essentially of, by weight, 60 to 75% SiO2, 17 to 27% Al2O3, 3 to 5% Li2O, 0 to 0.9% MgO, 0 to 0.9% ZnO, 0.3 to 5% BaO, 0 to 3% Na2O, 0 to 3% K2O, 0 to 4% TiO2, 1 to 4% ZrO2, 0 to 4% P2O5, 0.05 to 2% Sb2O3, 0 to 0.9% MgO+ZnO, and 0 to 4% Na2O+K2O."}
{"text": "1. Technical field\nThe present disclosure relates to components for a gas turbine engine and, more particularly, to the additive manufacture thereof.\n2. Background Information\nGas turbine engines typically include a compressor section to pressurize airflow, a combustor section to burn a hydrocarbon fuel in the presence of the pressurized air, and a turbine section to extract energy from the resultant combustion gases.\nIn the gas turbine industry, methods for fabricating components with internal passageways, such as blades and vanes within the turbine section, using additive manufacturing invite much attention. Since a component is produced in a continuous process in an additive manufacturing operation, features associated with conventional manufacturing processes such as machining, forging, welding, casting, etc. can be eliminated leading to savings in cost, material, and time.\nAn inherent feature of metallic components fabricated by additive manufacturing is that the metallic material forming the component has a polycrystalline microstructure. However, for numerous types of turbine components it is preferable to use a metallic material having a single crystal, or a columnar grain microstructure, which microstructure is able to withstand the higher temperatures and stresses typically experienced in the operating environment in a hot gas stream."}
{"text": "Typical hockey stick blades are generally made of a core reinforced with one or more layers of synthetic materials such as fiberglass, carbon fiber or Aramid. The core of the blade may also be made of a synthetic material reinforced with layers of fibers. The layers may be made of a woven filament fiber, preimpregnated with resin. Prior art structures have included a foam core with a layer of preimpregnated tape in the manner of pieces of bread on a sandwich, with a layer of tape on the top and bottom surfaces of the foam but not extending continuously along the edges of the core."}
{"text": "The administration of a therapeutic agent by injection (e.g., intravenous, intramuscular or subcutaneous injection) typically is regarded as the most convenient way of administration when the purpose is to achieve a rapid and strong systemic effect, for example, within 3-10 minutes, when the agent is not absorbed by the gastrointestinal tract, or when the agent is inactivated in the gastrointestinal tract or by first-pass hepatic metabolism. However, administration by injection presents a range of disadvantages. For example, sterile syringes must be used and injections cannot be administered by untrained personnel. Furthermore, this mode of administration may cause pain and/or irritation, especially in the case of repeated injections at the same site.\nMucosal administration, such as, intranasal, buccal, sublingual, rectal and pulmonal administration, is receiving particular interest as it avoids many of the disadvantages of injecting a therapeutic agent while, at the same time, still providing a strong and rapid systemic effect. In order to be an attractive alternative to injection, mucosal administration, for example, intranasal administration, should neither cause significant pain, discomfort or irritation nor cause any irreversible damage to the mucosal surface. However, in the case of acute health threatening indications, a relatively high local irritation to the mucosa may be acceptable.\nIn mucosal administration, such as during nasal, buccal or rectal administration, the therapeutic agent should be applied to the mucosa in a vehicle that permits it to penetrate, or be absorbed through, the mucosa. In order to penetrate the mucus, the vehicle should be biocompatible with mucus and hence have a certain degree of hydrophilicity. However, the vehicle should preferably also possess lipophilic properties to dissolve a clinically relevant amount of the therapeutic agent of interest.\nThe extensive network of blood capillaries under the mucosal surface, especially in the nasal mucosa, is well suited to provide a rapid and effective systemic absorption of drugs, vaccines and biologicals. Moreover, the nasal epithelial membrane in effect contains a single layer of epithelial cells (pseudostratified epithelium) and, therefore, is more suited for drug administration than other mucosal surfaces having squamous epithelial layers, such as, the mouth and vagina.\nIt has been hypothesized that the usefulness of nasal administration can be limited if the therapeutic agent has limited solubility in water (Proctor, D. F. (1985) Nasal Physiology in Intranasal Drug Administrations, in Chien, Y. W. (Ed.) TRANSNASAL SYSTEMIC MEDICATIONS, FUNDAMENTALS, DEVELOPMENTAL CONCEPTS AND BIOMEDICAL ASSESSMENTS, ELSEVIER Science Publishers, Amsterdam, pp. 101-105). As a result, this hypothesis, if correct, may limit the delivery of certain therapeutic agents that are sparingly soluble in water.\nTo facilitate delivery to the nasal cavity, an effective amount of the therapeutic agent should be dissolved in a small volume, for example, less than about 1000 μL, preferably less than 300 μL, and more preferably less than 150 μL. Larger volumes drain out anteriorly through the nostrils or posteriorly toward the pharynx where excess liquid is swallowed. As a result, if large volumes are administered, a portion of the therapeutic agent can be lost from the absorption site, and it can be difficult if not impossible to reproducibly administer the correct dose of the therapeutic agent.\nA variety of delivery systems have been developed for the nasal administration of therapeutic agents. Lau and Slattery studied the absorption characteristics of diazepam and lorazepam following their intranasal administration for the treatment of epilepticus (Lau, S. W. J. & Slattery, J. T. (1989), Absorption of Diazepam and Lorazepam Following Intranasal Administration, INT. J. PHARM., 54, 171-174). In order to solubilize the therapeutic agent, a non-ionic surfactant—polyoxyethylated castor oil—was selected as the least irritating solvent of several solvents studied, including polyethylene glycol 400 (PEG 400). Diazepam absorption was 84% and 72%, respectively, in two adults measured over a period of 60 hours. However, the peak concentration was not observed until 1.4 hours after the nasal administration and was only about 27% with reference to intravenous administration, suggesting that most of the absorption had taken place after the test substance passed down to pharynx and swallowed. Similar results were obtained for lorazepam but with an even longer time to peak (2.3 hours). The authors concluded that the intranasal route of administration had limited potential for the acute treatment of epileptic seizures.\nWilton et al. attempted to administer midazolam to 45 children to achieve pre-anesthetic sedation (Wilton et al. (1988) Preanaesthetic Sedation of Preschool Children Using Intranasal Midazolam, ANESTHESIOLOGY, 69, 972-975). However, the volumes used were impractical and exceeded the maximal volume required for efficient administration. This resulted in coughing and sneezing with expulsion of at least part of the dose.\nMorimoto et al. studied a gel preparation for nasal application in rats of nifedipine containing the gelling agent carbopol (polyacrylic acid) in PEG 400, for achieving prolonged action and high bioavailability of the therapeutic agent (Morimoto et al. (1987) Nasal Absorption of Nifedipine from Gel Preparations in Rats, CHEMICAL AND PHARMACEUTICAL BULLETINS, 35, No. 7, 3041-3044). A mixture of equal amounts of carbopol and PEG 400 was preferred. It was shown that nasal application provided higher bioavailability of nifedipine than after peroral administration, but the peak plasma concentration was not observed until 30 minutes after administration.\nDanish Patent Application No. 2586/87 discloses a pharmaceutical composition comprising an anti-inflammatory steroid, water, 2 to 10% (v/v) propylene glycol, 10 to 25% (v/v) PEG 400, and 1 to 4% (v/v) Tween 20.\nU.S. Pat. No. 4,153,689 discloses a stable aqueous solution of insulin intended for intranasal administration. The solutions had a pH not more than 4.7, and contained from 0.1 to 20% by weight of a stabilizing agent including (a) one or more non-ionic surface active agents whose hydrophile-lipophile balance value was in the range of 9 to 22, and/or (b) polyethylene glycol whose molecular weight was in the range of from 200 to 7500. Exemplary non-ionic surface active agents included polyoxyethylene fatty acid ester, a polyoxyethylene higher alcohol ether, a polyoxyethylene alkylphenyl ether, or a polyoxyethylene alkylphenyl ether, or a polyoxyethylene hydrogenated castor oil.\nInternational Patent Publication No. DK-2075/90 discloses the nasal administration of therapeutic agents, for example, benzodiazepines, in compositions containing n-glycofurol, a derivative of polyethyleneglycol, for mucosal administration. The application discloses the nasal administration of therapeutic agents, for example, benzodiazepines, in formulations containing at least 30% n-ethyleneglycols ranging from 1-8 ethylene glycol, for example, polyethylene glycol 200 (PEG 200).\nU.S. Pat. No. 5,693,608 discloses a method of administering a therapeutic agent via the nasal mucosa of a mammal, where the agent is dissolved or suspended in an n-ethyleneglycol containing vehicle where the n-ethyleneglycol is represented by the formula, H(OCH2CH2)pOH, wherein p is a number from 1 to 8.\nNotwithstanding, there is still a need for compositions deliverable through mucosal membranes that produce therapeutic plasma concentrations of the therapeutic agent as fast as or nearly as fast as by intravenous administration but without causing irritation and/or unacceptable damage to the mucosal membrane."}
{"text": "1. Field of the Invention\nThe present invention relates to a point-of-sale system for retail environments, and more particularly to a system and method for detecting and analyzing motion patterns of individuals at a self-service checkout point-of-sale.\n2. Description of the Prior Art\nIn the retail industry, labor is the number one expense after cost of goods sold. Employees working as \"checkout personnel\" or \"cashiers\" are the single largest category of labor. This is particularly true for those retailers in the food segment. On average, the retailer must have three or more cashiers trained for every checkout lane in a food store.\nOne problem for retailers is the cost associated with the recruiting, hiring, training, scheduling, managing, etc. of cashiers. This cost is exacerbated in that turnover in the food segment can easily exceed 100% annually, especially in metropolitan areas. It is not uncommon for food retailers to have significantly fewer cashiers on hand in a retail establishment than is optimally required.\nThe constant recruiting, hiring, training, and scheduling of cashiers consumers significant management time and expense. Scheduling alone can make or break a retailer's performance on any given day and significantly impact customer service. By scheduling too many cashiers, the retailer has excess capacity and a higher cost of sales for the day. By scheduling too few cashiers, checkout queues grow long with angry customers who may not return due to poor customer service. Other customers may enter the store and go elsewhere due to the long lines. Checkout problems are typically the leading consumer dissatisfaction issue for retailers.\nSelf Checkout (SCO) Terminals offer benefits to both the retailer and the consumer. For the retailer, SCO terminals reduce the retailer's need for cashiers, thereby reducing the costs associated with them. Since SCO terminals are rarely closed, the job of predicting and scheduling cashier demand is made easier. For the consumer, SCO terminals offer the perception of faster checkout, privacy, control, etc., and a significantly enhanced shopping experience.\nThus, in theory, replacing traditional cashier assisted checkout with self checkout makes great business sense. However, the cashier performs several roles for the retailer in addition to itemizing and tendering groceries. One significant role is that of security--ensuring that the consumers pay for all the groceries that they take out of the store. A SCO terminal can enable the consumer to take the cashier's place in itemizing and tendering. But the concern remains how to ensure that the consumer pays for all the items taken out of the store.\nLoss prevention specialists know that the probability of theft increases when someone believes an opportunity exists to deceive without an offsetting risk of detection/consequence. By removing the cashier, the consumer is no longer faced with a significant method of detection.\nTypical SCO solutions have used technologies that are not consumer friendly, are perceived as non-trusting, fail to support the consumer throughout the transaction, and have placed added operational issues on the retailer.\nThere is therefore a significant need in the art for a system and method for providing security for self checkout environments in a way that is effective and efficient, but is also acceptable to the consumer."}
{"text": "The invention relates to a method for order picking products in an order picking facility having at least one rack for storing the products and an assigned warehouse transport system for collecting and storing the products and for transferring the collected products to a transport system, by means of which the collected products are transported to a good-to-order picking station, at which the quantity of the collected products that is determined in a picking order is manually or automatically by a robot, respectively, order-picked into one order transport unit for each picking order, wherein product types corresponding to the picking order are collected, order-picked at the good-to-picker order picking station (WZM1) one after the other and wherein the remaining products of the product type are re-stored via the transport system again in the rack.\nThe invention further relates to an order picking facility for order picking products, wherein the order picking facility has at least one rack for storing the products and a warehouse transport system assigned to the rack for collecting and storing and the products and for transferring the collected products to a transport system and wherein there is provided a warehouse administration computer for controlling the order picking facility, and wherein there is provided a good-to-picker order picking station, at which the quantity of the collected goods that is determined in a picking order may be order picked in one order transport unit for each picking order, wherein product types corresponding to the picking order may be collected by the warehouse transport system, order picked at the good-to-picker order picking station and wherein products of the product type remaining after the order picking may be re-stored via the transport system in the rack.\nSuch a method and such an order picking facility are, for example, known from the document DE 101 36 354 A1. The known system comprises several high racks of a high-rack storage, between which there is formed a lane width, in which there is provided respectively one warehouse transport system in order to store or collect the products stored in the high-rack storage. Each of the warehouse transport systems has one level operating tool per rack level, by means of which the products to be collected are picked up at storage sites and transferred to a lift of the transport system. By means of conveyor belts of the transport system, collected products are conveyed to the good-to-picker order picking station, at which an operator takes the quantity of the product type of the collected product that is determined for a picking order and places this into an order container or an order transport unit, respectively, of the picking order. A warehouse administration computer of the order picking facility controls the warehouse transport systems and the transport system so that all products that are to be order picked for a picking order are collected from the high-rack storage one after the other and transported to the good-to-picker order picking station. The operator then only has to take out the quantity of the product type of the collected products that is displayed on a monitor or by light symbols and put these into the order container A. Collected products remaining after the order picking will then be re-stored in the high-rack storage, controlled by the warehouse administration computer.\nIn the known method for order picking and the known order picking facility it has been proven that this way of good-to-picker order picking may be rather satisfyingly applied for certain picking orders. It is, however, disadvantageous that for certain product types or certain product characteristics (e.g., turn-over frequency, product geometry), respectively, or in order to obtain high order picking through-put performances for certain picking orders or parts of certain picking orders there may be applied order picking systems that are more suitable than the good-to-picker order picking system. Order picking by means of central belt technology, for example, has been rather successfully used for the order picking of, e.g., medicinal products. In central belt technology, there is provided a completely different storage system according to the state of art, which is why in known order picking facilities there has to be decided immediately after the delivery of the products by truck or rail, how the products are to be finally order picked in order to determine the suitable storage type in the order picking facility."}
{"text": "1. Field of the Invention\nThis invention relates to the electrostatic spraying of liquids in such a way that the liquid is initially projected from a spray head in the form of a ligament which thereafter breaks up into droplets under the influence of Coulombic forces to produce an atomised spray. Electrostatic spraying of this type is well known and is described in for example our prior British Patent No. 1569707.\n2. Description of the Related Art\nIn conventional ligament mode spraying, it is widely recognised that liquid resistivity is vitally important to securing satisfactory atomisation and that aqueous and other liquids which have relatively low resistivities become more and more unsuitable for use in ligament mode spraying as resistivity reduces below 1.times.10.sup.7 ohm cm.\nAlthough not limited thereto, the present invention is particularly concerned with the spraying of relatively low resistivity liquids such as aqueous, alcohol and aqueous/alcohol based liquids commonly used in personal care products such as deodorants, anti-perspirants, scents and hair sprays. In the past, many such products have been marketed as aerosol products in which a propellant is used to cause atomisation of the liquid into fine droplets typically less than 50 microns in diameter.\nHowever, because of the currently perceived environmental problems associated with the propellants conventionally used in aerosols, attention has turned to alternative methods of dispensing personal care liquids. Electrostatic spraying offers one alternative approach, but, where the ingredient to be dispensed is combined with an aqueous and/or alcohol carrier (or other relatively low resistivity liquid), current wisdom suggests that, with practical flow rates (typically several cc/min), such carriers will not allow dispensing of the product as droplets with a size range comparable to that attainable with aerosol sprays.\nEP-A-152446 discloses a device for the electrostatic spraying of aqueous liquids and explains that, for reasons not completely understood, satisfactory atomisation of aqueous formulations can only be achieved at flow rates that are undesirably low for many purposes and ligamentary formation is not obtained with aqueous liquids. EP-A-152446 proposes the use of a corona discharge needle electrode assembly in the vicinity of a sprayhead including a narrow metal tube having a diameter of 400 microns, the arrangement being such that the electrode assembly is symmetrically disposed about the emerging liquid and produces ions which bombard the liquid so that the liquid assumes a stable ligamentary form. It is stated that the illustrated embodiment produces droplets having a volume median diameter of 10 to 50 microns. For personal care products and like products for domestic use, it is considered undesirable to locate an assembly of needle electrodes in the vicinity of the outlet of the device both from an aesthetic standpoint and also in terms of the risk of potential electrostatic shock."}
{"text": "1. Field of the Invention\nEmbodiments of the invention relate to electronic data transmission, and more particularly, in one or more embodiments, to multi-level signaling.\n2. Description of the Related Art\nIn electronic data transmission, signals representing electronic data are transmitted from a transmitting component to a receiving component. The electronic data typically includes data symbols such as binary digits (which are often referred to as bits), i.e., 0's and 1's. In many applications, signals having two different voltage levels representing data symbols such as 0's and 1's have been widely used for such data transmission.\nIn certain applications, signals having more than two levels have been used to transmit electronic data. Such a signaling scheme is generally referred to as “multi-level signaling.” In a multi-level signaling scheme, the data rate can be increased without increasing the clock frequency or the number of channels. In such a scheme, transmitters generate multi-level signals, and receivers detect them, allowing multiple (k) bits to be transmitted or received as one of 2k possible voltages at each clock edge or at least once per clock cycle. A multi-level signaling scheme allows two or more bits of data to be simultaneously transmitted by multiple signal levels, thereby increasing the data throughput for a given operating frequency.\nIn certain instances, the term “multi-level signaling” is also referred to as multi-level pulse amplitude modulation or M-PAM signaling. In such a case, the number of signal levels is indicated by a number before the acronym “PAM.” For example, a PAM signaling scheme using four signal levels is represented by 4 PAM. Similarly, a PAM signaling scheme using eight signal levels is represented by 8 PAM.\nFIG. 1 is a graph illustrating a conventional multi-level signaling scheme using four voltage levels, i.e., 4 PAM. The four signal levels represent two-bits b0 b1. The highest voltage level over a high reference voltage VREFH represents the bits “11.” The second highest voltage level between the high reference voltage VREFH and a middle reference voltage VREFM (which is lower than VREFH) represents the bits “10.” The third highest voltage level between the middle reference voltage and a low reference voltage VREFL (which is lower than VREFM) represents the bits “01.” The lowest voltage level lower than VREFL represents the bits “00.” The two bits are transmitted as a single multi-level symbol at every clock edge by transferring an appropriate one of the four voltage levels. Therefore, the data rate of the signaling scheme just described is twice that of a 2-PAM system.\nReferring to FIG. 2, a conventional transmitter circuit 200 for generating voltage levels under a 4 PAM signaling scheme will now be described. To provide the voltage levels to transmit a 4-PAM symbol, the transmitter circuit 200 sinks a predetermined amount of current for that symbol. In particular, each symbol is associated with a distinct amount of current.\nThe transmitter circuit 200 includes a voltage source VDD, a resistor R, a first transistor TR1, a second transistor TR2, a first current source CS1, a second current source CS2, a first node N1, and a signal output VOUT. The resistor R is connected between the voltage source VDD and the first node N1. The first node N1 is electrically connected to the signal output VOUT. Each of the first and second transistors TR1, TR2 is connected to the first node N1 at its source/drain. The drain/source of the first transistor TR1 is connected to the first current source CR1 which provides a current of 2I. The drain/source of the second transistor TR2 is connected to the second current source CR2 which provides a current of I.\nTo transmit the bits “11, ” the transmitter circuit 200 sinks no current by turning off both of the first and second transistors TR1, TR2, and the signal output VOUT is pulled up to VDD. To transmit the bits “10, ” the transmitter circuit 200 sinks an amount of current I by turning on the second transistor TR2 only, thereby providing VDD−RI at the signal output VOUT. To transmit the bits “01, ” the transmitter circuit 200 sinks an amount of current 2I by turning the first transistor TR1 only, thereby providing VDD−2RI at the signal output VOUT. To transmit the bits “00, ” the transmitter circuit 200 sinks an amount of current 3I by turning on both of the first and second transistors TR1, TR2, thereby providing VDD−3RI at the signal output VOUT.\nIn the transmitter described above, when generating three out of the four signal levels (e.g., those representing “00”, “01”, and “10”), power is dissipated because there is a current flow through the resistor R in the circuit. In other words, three out of four signal levels consume static power. It should be understood that the translation of the two sequential binary bits into multiple voltage levels need not follow exactly as it has been described thus far. For example, the bits “00” could be encoded as the highest voltage level, rather than as the lowest level, as has thus far been illustrated. Other alternative exist as well, including the well known gray-coding which would order the levels as follows: “00”, “01”, “11”, and “10.”"}
{"text": "FIG. 1 shows a configuration of a known door 10 for the vehicle. A door frame 1 forms as an edge of a window portion 4 of the door 10. A cross sectional view of a known common door frame 1 along a line A—A in FIG. 1 is shown in FIG. 5 as an example. Specifically, the door frame 1 made by bending a plate member comprises a frame structure portion 5 with a hollow portion 51 for being provided with its stiffness and a flange portion 54 extending from the outer portion in one direction, and an outer surface portion 6 formed on the top portion of the flange portion 54 in FIG. 5. The outer surface portion 6 is formed in approximately flat shape and assembled to be viewed from outside of the vehicle. A glass run groove 52 is formed on the right of the flange portion 54 in horizontal direction in FIG. 5 whereinto a glass run 31 inserted for engaging with a window glass 3. A seal groove 53 is formed on the left of the flange portion 54 for supporting a seal material 8 between the door 10 and the vehicle body. In this configuration, the flange portion 54 includes three layered palates which are welded each other at a welded place 9 positioned at an approximately center portion of the flange portion 54 as shown in FIG. 5.\nAccording to the door frame 1 which is generally configured as described above, many kinds of inventions have been reported by which the weight of the door frame 1 is saved. One example of these fames comprises individually formed frame structure portion 5 and outer surface portion 6, and they are welded each other, so that each sizes of the hollow portion 51 and the outer surface portion 6 in cross sectional form are adapted to be changed in longitudinal direction. (e.g. Tokukaihei 9-216583 kouhou)\nIn addition, another example of these frames is formed by a sheet of plate member which is made of welded plates with difference thickness, and bended in aforementioned form. (e.g. Tokukaihei 10-297280 kouhou)\nHowever, in the known aforementioned frame of the vehicle door, a flange portion comprising three-layered plates, which result in increasing its weight, is not necessarily for ensuring the strength of the frame. In addition, a distortion of the flange portion generated by the welding causes a distortion on the outer surface portion; as a result, it becomes difficult to achieve a smooth outer surface portion with a quality required for vehicle appearance.\nAccordingly, it is an object of the present invention to provide an improved door frame structure for ensuring a quality of the outer surface portion and weight saving in which the outer surface portion is not affected by the distortion."}
{"text": "As noted in U.S. Pat. No. 3,985,324, issued Oct. 12, 1976 to Larson, garden tools such as rakes, shovels, hand hoes and the like are often stored in locations where the tools may become a hazard to persons in their vicinity. Unwary or careless people may step on a sharp or pointed tool end, or brush against a precariously balanced implement. These encounters invariably cause aggravation, and can cause serious injury leading to otherwise avoidable medical bills and risk to life, limb and property.\nTo eliminate such hazards, many people hang such tools on hooks, pegs, nails, or other wall mounts. However, the amount of wall surface available for hanging such tools is often limited, and it is common practice to store a number of tools in stacked relationship upon the same hook or peg. Stacked and hung tools are also hazardous since the endmost tool may fall from the stack, and since the supportive mount may be unable to carry the aggregate weight of these tools. Thus, there is a need for new and/or improved tool holders."}
{"text": "The present invention relates to implantable medical devices (IMDs). Specifically, the invention pertains to a remote bi-directional communications between the IMDs and a drug dispenser. More specifically the invention pertains to a closed loop system in which the IMDs monitor and determine the presence of a specific drug dose in the patient\"\"s body to send instructions to the drug dispenser or an interface medical unit (IMU) to implement a drug management scheme based on the monitored data. More specifically, the invention provides a dynamic drug management system in which the drug dose is chronically monitored by the IMDs to enhance drug effectiveness and as well monitor patient compliance with recommended drug administration regimen. The invention preferably utilizes a robust communication system integrated with a remote expert data center in a web-enabled environment to transmit the IMDs\"\" data to a physician for evaluation and review thereby enhancing the delivery of therapy and clinical care remotely.\nA technology-based health care system that fully integrates the technical and social aspects of patient care and therapy should be able to flawlessly connect the client with care providers irrespective of separation distance or location of the participants. While clinicians will continue to treat patients in accordance with accepted modern medical practice, developments in communications technology are making it ever more possible to provide a seamless system of remote patient diagnostics, care and medical services in a time and place independent manner.\nPrior art methods of clinical services are generally limited to in-hospital operations. For example, if a physician needs to review the performance parameters of an implantable device in a patient, it is likely that the patient has to go to the clinic. Further, if the medical conditions of a patient with an implantable device warrant a continuous monitoring or adjustment of the device, the patient would have to stay in a hospital indefinitely. Further, if the patient with the IMDs is taking a drug, it is often clinically prudent to monitor the dose and its impact on the patient and, as well, on the IMDs. Such a continued treatment plan poses both economic and social problems. Under the exemplary scenario, as the segment of the population with implanted medical devices increases many more hospitals/clinics including service personnel will be needed to provide in-hospital service for the patients, thus escalating the cost of healthcare. Additionally the patients will be unduly restricted and inconvenienced by the need to either stay in the hospital or make very frequent visits to a clinic.\nYet another condition of the prior art practice requires that a patient visit a clinic center for occasional retrieval of data from the implanted device to assess the operations of the device and gather patient history for both clinical and research purposes. Such data is acquired by having the patient in a hospital/clinic to down load the stored data from the implantable medical device. Depending on the frequency of data collection this procedure may pose serious difficulty and inconvenience for patients who live in rural areas or have limited mobility. Similarly, in the event a need arises to upgrade the software of an implantable medical device, the patient will be required to come into the clinic or hospital to have the upgrade installed. Further, in medical practice it is an industry-wide standard to keep an accurate record of past and contemporaneous procedures relating to an IMD uplink with, for example, a programmer. It is required that the report contain the identification of all the medical devices involved in any interactive procedure. Specifically, all peripheral and major devices that are used in down linking to the IMD need to be reported. Currently, such procedures are manually reported and require an operator or a medical person to diligently enter data during each procedure. One of the limitations of the problems with the reporting procedures is the fact that it is error prone and requires rechecking of the data to verify accuracy.\nA further limitation of the prior art relates to the management of multiple medical devices in a single patient. Advances in modern patient therapy and treatment have made it possible to implant a number of devices in a patient. For example, IMDs such as a defibrillator or a pacer, a neural implant, a drug pump, a separate physiologic monitor and various other IMDs may be implanted in a single patient. To successfully manage the operations and assess the performance of each device in a patient with multi-implants requires a continuous update and monitoring of the devices. As is often the case, patients with multi-implanted medical devices may take a variety of medications. It is therefore necessary to monitor drug intake and its effect on the oprerational and functional parameters of the IMDs. More importantly, chronic monitoring of drug intake and its effect on the physiological and clinical conditions of the patient enables a proactive intervention to change the course of an otherwise serious medical condition. Thus, there is a need to monitor drug delivery and effectiveness via IMDs.\nAccordingly it is vital to have a drug dispenser unit that would establish a communication system with IMDs. The unique position of IMDs enables a real-time assessment of physiological conditions which may change or indicate a measurable variance due to drug dose and delivery. IMDs could be adapted to provide measurements relating to the physiological impact of drug therapy. Further, IMDs could be adapted to provide a quick evaluation of the effectiveness of a drug to support a clinical decision as to whether a given dose is a prudent course of therapy.\nThe proliferation of patients with multi-implant medical devices worldwide has made it imperative to provide remote services to the IMDs and timely clinical care to the patient. Frequent use of programmers to communicate with the IMDs and provide various remote services, consistent with co-pending applications titled xe2x80x9cApparatus and Method for Remote Troubleshooting, Maintenance and Upgrade of Implantable Device Systems,xe2x80x9d filed on Oct. 26, 1999, Ser. No. 09/426,741; xe2x80x9cTactile Feedback for Indicating Validity of Communication Link with an Implantable Medical Device,xe2x80x9d filed Oct. 29, 1999, Ser. No. 09/430,708; xe2x80x9cApparatus and Method for Automated Invoicing of Medical Device Systems,xe2x80x9d filed Oct. 29, 1999, Ser. No. 09/430,208; xe2x80x9cApparatus and Method for Remote Self-Identification of Components in Medical Device Systems,xe2x80x9d filed Oct. 29, 1999, Ser. No. 09/429,956; xe2x80x9cApparatus and Method to Automate Remote Software Updates of Medical Device Systems,xe2x80x9d filed Oct. 29, 1999, Ser. No. 09/429,960; xe2x80x9cMethod and Apparatus to Secure Data Transfer From Medical Device Systems,xe2x80x9d filed Nov. 2, 1999, Ser. No. 09/431,881; xe2x80x9cImplantable Medical Device Programming Apparatus Having An Auxiliary Component Storage Compartment,xe2x80x9d filed Nov. 4, 1999, Ser. No. 09/433,477; which are all incorporated by reference herein in their entirety, has become an important aspect of patient care. Thus, in light of the referenced disclosures, communication with IMDs enhances the delivery of therapy and clinical care in real time. Specifically, as the number of patients with IMDs increases globally, the need to manage drug delivery and intake remotely becomes an economic imperative. Further, IMDs which are communicable and operable in a web-enabled environment, as contemplated by the cited disclosures hereinabove, provide a unique platform to assess the efficacy of drugs and the compliance of patients with prescribed regimens. Further, it is vital to have a drug dispenser that is adapted to have data communications with the IMDs and other data centers to support the remote patient management system contemplated by the present invention.\nThe prior art provides various types of remote sensing and communications with an implanted medical device. One such system is, for example, disclosed in Funke, U.S. Pat. No. 4,987,897 issued Jan. 29, 1991. This patent discloses a system that is at least partially implanted into a living body with a minimum of two implanted devices interconnected by a communication transmission channel. The invention further discloses wireless communications between an external medical device/programmer and the implanted devices.\nOne of the limitations of the system disclosed in the Funke patent includes the lack of communication between the implanted devices, including the programmer, with a remote clinical station. If, for example, any assessment, monitoring or maintenance is required to be performed on the IMD the patient will have to go to the remote clinic station or the programmer device needs to be brought to the patient\"\"s location. More significantly, the operational worthiness and integrity of the programmer cannot be evaluated remotely thus making it unreliable over time as it interacts with the IMD. Further, in light of the present invetion, the Funke patent does neither suggest nor disclose the communications system between the IMD and a drug dispenser to monitor and assess in the effectiveness of the dose based on the physiological status of the patient.\nYet another example of drug management based on smart drug dispenser units is disclosed by Martindale et al in U.S. Pat. No. 4,360,125 issued on Nov. 23, 1982. In the disclosure, a medication dispenser in which medication to be dispensed is housed including a member operable to allow medication access. The dispenser provides a medication alert signal at preselected times in accordance with a desired medication regimen. A medication access signal is provided when medication access is obtained. Data representative of the relative timing between a medication alert signal and a medication access signal is written into readable memory whereby that data is available to a physician for evaluation. In the preferred embodiment, the data is representative of the time of occurrence of each medication alert signal and medication access signal. The interval between medication alert signals is selectively alterable.\nFurther, examples of drug management based on smart drug dispensers are disclosed in U.S. Pat. Nos. 4,768,176; 4,768,177; 5,200,891; 5,642,731; 5,752,235 and 5,954,641 all to Kehr et al. Generally all the patents relate to a drug dispensing system with various alert features to monitor and manage the administration of medication and medical treatment regimens. None of these patents suggest or disclose a communication between the drug dispensing systems and an IMD.\nYet another prior art reference provides a multi-module medication delivery system as disclosed by Fischell in U.S. Pat. No. 4,494,950 issued Jan. 22, 1985. The disclosure relates to a system consisting a multiplicity of separate modules that collectively perform a useful biomedical purpose. The modules communicate with each other without the use of interconnecting wires. All the modules may be installed intracorporeal or mounted extracorporeal to the patient. In the alternate, some modules may be intracorporeal with others being extracorporeal. Signals are sent from one module to the other by electromagnetic waves. Physiologic sensor measurements sent from a first module cause a second module to perform some function in a closed loop manner. One extracorporeal module can provide electrical power to an intracorporeal module to operate a data transfer unit for transferring data to the external module.\nThe Fischell disclosure provides modular communication and cooperation between various medication delivery systems. However, the disclosure does not provide an external pill dispenser which is in wireless communications with IMDs. Further, the system does neither teach nor disclose an external programmer for telemetrically interacting with the pill dispenser.\nAccordingly, it would be advantageous to provide a pill dispenser that communicates with IMDs to implement an effective drug management system. Yet another desirable advantage would be to provide a high speed communications scheme to enable the transmission of high fidelity sound, video and data to advance and implement efficient remote drug management of a clinical/therapy system via a programmer thereby enhancing patient clinical care. As discussed herein below, the present invention provides these and other desirable advantages.\nThe present invention generally relates to a communications scheme in which a remote web-based expert data center interacts with a patient having one or more implantable medical devices (IMDs) via an associated external medical device, preferably a programmer, located in close proximity to the IMDs. The IMDs are adapted to communicate with a pill dispenser to monitor and log pill deposition and effectiveness. Some of the most significant advantages of the invention include the use of various communications media between the remote web-based expert data center and the programmer to remotely exchange clinically significant information and ultimately effect real-time drug intake and prescriptive changes as needed.\nOne of the many aspects of the present invention includes a real-time access of a programmer or a pill dispenser to a remote web-based expert data center, via a communication network, which includes the Internet. The operative structure of the invention includes the remote web-based expert data center, in which an expert system is maintained, having a bi-directional real-time data, sound and video communications with the programmer via a broad range of communication link systems. The programmer is in turn in telemetric communications with the IMDs such that the IMDs may uplink to the programmer or the programmer may down link to the IMDs, as needed.\nYet another feature of the invention includes a communications scheme that provides a highly integrated and efficient method and structure of clinical information management in which various networks such as Community access Television, Local area Network (LAN), a wide area network (WAN) Integrated Services Digital Network (ISDN), the Public Switched telephone Network (PSTN), the Internet, a wireless network, an asynchronous transfer mode (ATM) network, a laser wave network, satellite, mobile and other similar networks are implemented to transfer voice, data and video between the remote data center and a programmer. In the preferred embodiment, wireless communications systems, a modem and laser wave systems are illustrated as examples only and should be viewed without limiting the invention to these types of communications alone. Further, in the interest of simplicity, the applicants refer to the various communications system, in relevant parts, as a communications system. However, it should be noted that the communication systems, in the context of this invention, are interchangeable and may relate to various schemes of cable, fiber optics, microwave, radio, laser and similar communications or any practical combinations thereof.\nSome of the distinguishing features of the present invention include the use of a robust web-based expert data center to collect drug therapy information based on data communication between the IMDs, the pill dispenser and the programmer. Specifcally the invention enables remote evaluation of drug performance in a patient. Although the present invention focuses on the remote real-time monitoring and management of drug therapy information, the system could advantageously be used to monitor clinical trials of drugs or collect clinical data relating to drug interaction or physiological impact of various doses on the patient.\nYet one of the other distinguishing features of the invention includes the use a highly flexible and adaptable communications scheme to promote continuous and real-time communications between a remote expert data center, a programmer and a pill dispenser associated with a plurality of IMDs. The IMDs are structured to share information intracorporeally and may interact with the programmer or the pill dispenser as a unit. Specifically, the IMDs either jointly or severally can be interrogated to implement or extract clinical information as required. In other words, all of the IMDs may be accessed via one IMD or, in the alternate, each one of the IMDs may be accessed individually. The information collected in this manner may be transferred to the data center via the programmer or pill dispenser by up linking the IMDs as needed.\nThe invention provides significant compatibility and scalability to other web-based applications such as telemedicine and emerging web-based technologies such as tele-immersion. For example, the system may be adapted to webtop applications in which a webtop unit may be used to uplink the patient to a remote data center for drug information exchange between the IMDs and the remote expert data center. In these and other web-based similar applications the data collected, in the manner and substance of the present invention, may be used as a preliminary screening to identify the need for further intervention using the advanced web technologies.\nMore significantly, the invention provides a system and method to remotely monitor drug effectiveness in a patient. Further, the invention enables a chronic evaluation of drugs in a patient on real time basis. The significance of this method includes the fact that the data collected in this manner could be used to influence the course of drug therapy."}
{"text": "(a) Field of the invention:\nThe present invention relates to a large-scale semiconductor integrated circuit device (called LSI hereinafter) and more particularly to an LSI made by a standard-cell system.\n(b) Description of the prior art:\nIn order to construct various hardware systems, there have been already used LSI's of various types shown in FIG. 1, that is, standard logics such as TTL, microprocessor, semi-custom LSI and full-custom LSI. Among them, the types which can realize logic designs as desired are a semi-custom LSI including ROM, PLA (Programmable Logic Array), gate array and standard-cell LSI, and a full-custom LSI of manual design. The gate array LSI is manufactured to change only the final metallization steps of the manufacturing process using a master slice, thereby realizing the desired logic ; and the full-custom LSI is manufactured in accordance with any desired specification from the first step of the manufacturing process.\nOf the semi-custom LSI, the standard-cell LSI is designed by selecting proper standard cells from a standard cell library which is stored in a computer and in which are registered in advance such various standard cells having any desired logical functions as counter, buffer, etc. comprising plural circuit elements such as transistor, resistor, etc. An example of the standard cell library is shown in Table 1. The layout on a semiconductor chip is conducted, by means of a CAD/CAE (Computer Aided Design/ Computer Aided Engineering) system, so that the standard cells are automatically placed at proper positions on the semiconductor chip and the routing between the respective standard cells is also automatically carried out. Therefore, the standard-cell LSI is a semi-custom LSI standardized in its constituent elements but is manufactured according to any desired specification from the first step of the manufacturing process as in the full-custom LSI in the manufacture.\nTABLE 1 __________________________________________________________________________ (STANDRAD CELL LIBRARY OF SITL LSI) GROUP FUNCTION NAME __________________________________________________________________________ GATE INVERTER, NAND 1 OUTPUT NOT1 INVERTER, NAND 2 OUTPUT NOT2 INVERTER, NAND 3 OUTPUT NOT3 INVERTER, NAND 4 OUTPUT NOT4 2 INPUT OR OR2 3 INPUT OR OR3 4 INPUT OR OR4 2 INPUT NOR NOR2 3 INPUT NOR NOR3 4 INPUT NOR NOR4 EXCLUSIVE OR EOR EXCLUSIVE NOR EORN BUFFER POWER INVERTER NOTP (POWER GATE) NON INVERTING BUFFER 1 OUTPUT BUF1 NON INVERTING BUFFER 2 OUTPUT BUF2 MULTIPLEXER 8 TO 1 MULTIPLEXER MPX1 4 TO 1 MULTIPLEXER MPX2 DUAL 4 TO 1 MULTIPLEXER MPX3 2 TO 1 MULTIPLEXER MPX4 QUAD 2 TO 1 MULTIPLEXER MPX5 FLIP-FLOP D FLIP-FLOP DF1 D FLIP-FLOP WITH OE DF1E D FLIP-FLOP WITH RESET DFR D FLIP-FLOP WITH RESET & OE DFRE D FLIP-FLOP WITH SET & RESET DFQ D FLIP-FLOP WITH SET, RESET & OE DFQE JK FLIP-FLOP JK JK FLIP-FLOP WITH OE JKE SHIFT- 1 BIT SYNC LOAD SR SPL REGISTER 1 BIT ASYNC LOAD SR SP LATCH TRANSPARENT LATCH LC TRANSPARENT LATCH WITH OE LCE TRANSPARENT LATCH WITH RESET LCR TRANSPARENT LATCH WITH RESET & OE LCRE COUNTER 4 BIT SYNC UP COUNTER CTUL 4 BIT SYNC UP COUNTER CTU WITH DIRECT CLEAR 4 BIT SYNC UP/DOWN COUNTER CTUD 4 BIT SYNC UP COUNTER CTR WITHOUT PRESET 1 BIT RIPPLE COUNTER TFQ ADDER 1 BIT FULL ADDER FA 4 BIT FULL ADDER WITH FAST CARRY FA4 1 BIT HALF ADDER HA COMPARATOR 4 BIT MAGNITUDE COMPARATOR CMP DECODER 3 TO 8 DECODER DEC1 DECODER 2 TO 4 DECODER DEC2 DELAY 8 GATES DELAY DLY8 ELEMENT I/O BUFFER LS INPUT BUFFER IBLS S INPUT BUFFER IBS LS OUTPUT BUFFER OBLS S OUTPUT BUFFER OBS LS I/O BUFFER IOBLS LS TRI-STATE OUTPUT BUFFER OBLST LS SCHMITT TRIGGER INPUT BUFFER SIB LS LINE DRIVER OBLSL LS OPEN DRAIN OUTPUT BUFFER OBLSO TRI-STATE CONTROL BUFFER BUFTC SPECIAL CELL CRYSTAL OSCILLATOR XOS ROM 32 WORD .times. 8 Bit (256 Bit) ROM0 128 WORD .times. 8 Bit (1 kBit) ROM1 512 WORD .times. 8 Bit (4 kBit) ROM4 RAM 32 WORD .times. 8 Bit (256 Bit) RAM1 __________________________________________________________________________\nHere, the manual design full-custom LSI, standard-cell LSI and gate array are compared in respect of a plurality of items in Table 2.\nTABLE 2 __________________________________________________________________________ FULL-CUSTOM SEMI-CUSTOM ITEM MANUAL STANDARD-CELL GATE-ARRAY __________________________________________________________________________ LSI DESIGN BY USERS x .circle. .circle. UNNECESSITY OF x .circle. .circle. KNOWLEDGE OF PROCESS AND DEVICE FREEDOM OF FUNCTION .circleincircle. .circle. x (MEMORY, etc.) LAYOUT ERROR x .circle. .circle. LAYOUT COST x .circle. .circle. LAYOUT PERIOD x .circle. .circle. FREEDOM OF .circle. .circle. x GATE NUMBER OPTIMIZATION OF .circle. .circle. x I/O BUFFER (SPEEED, DRIVING POWER) FREEDOM OF .circle. .circle. x CHARACTERISTICS (SCHMITT TRIGGER, etc. DEVELOPMENT COST x .circle. .circle. DEVELOPMENT TIME x .circle. .circle. INTEGRATION DEGREE .circleincircle. .circle. x (CHIP COST) SPEED CHARACTERISTIC .circleincircle. .circle. x KEEPING OF .circle. .circleincircle. .circleincircle. CONFIDENTIALITY MULTI-FUNCTION .circle. .circle. x __________________________________________________________________________ x NOT GOOD, .circle. GOOD, VERY GOOD\nAccording to Table 2, the standard-cell LSI has following advantages over the gate array. The gate array is formed of a logical block called a macro cell which is constituted by a combination of general-purpose basic gates and which therefore contains gates which are eventually not used; on the other hand, the standard-cell LSI is formed of standard cells having respectively only specific functions, and has no gate which is not used, that is, no excess gate. Therefore, for example, in the case of constructing a same flip-flop, the standard cell LSI will have an occupied area smaller than that of the gate array. Also, it is possible to provide as standard cells CPU, RAM, ROM, PLA, etc., resulting in large freedom of function. Further, in the standard-cell system, there is no restriction on number of gates and the chip can be comprised of properly requisite number of gates. Various speeds and current driving power characteristics are also possible and a buffer such as a Schmitt trigger can be available in the standard-cell system. By the way, the PLA is not shown in Table 2 but is greatly restricted in the integration degree and function, and is therefore found to be inadequate to be made large on the scale.\nThe standard-cell LSI has many advantages as described above but, at the same time, has a defect that the logic and timing of the first designed specification can not be easily changed. It is recently so increasingly required to develop devices within a short turn around time. To this end, it is required to be able to first develop an LSI with a provisional specification and to finally alter the specification such as logic, timing, etc. of the LSI. The conventional standard-cell LSI has not been able to satisfy such requirement. In the conventional standard-cell LSI, the LSI must be manufactured newly from the first step of the manufacturing process after altering the design. Thus, in the conventional standard-cell LSI, there has been a great problem that, in order to change the specification such as logic, timing, etc., the development time and cost will increase."}
{"text": "Recent developments in the automotive industry have resulted in an increase in the demands on the properties of power transmission belts. As a result, power transmission belts on these engines have been required to operate under higher loads, at higher speeds and at higher temperatures while reducing belt width and noise. The automotive industry demands a high quality belt capable of withstanding these severe conditions. One method to improve the properties of such belts has been to improve the materials from which the belts are made. For example, conventional polychloroprene rubber has been replaced with hydrogenated nitrile rubber and chlorosulfonated polyethylene. In addition, solid lubricants have been added to the rubber compositions to improve the abrasion properties of the belt by reducing the coefficient of friction.\nThe present invention relates to a power transmission belt that is characterized by a layer of thermoplastic located between a fabric layer on the drive surface and the cushion section of the belt."}
{"text": "The complement system plays a central role in the clearance of immune complexes and the immune response to infectious agents, foreign antigens, virus-infected cells and tumor cells. However, complement is also involved in pathological inflammation and in autoimmune diseases. Therefore, inhibition of excessive or uncontrolled activation of the complement cascade could provide clinical benefit to patients with such diseases and conditions.\nThe complement system encompasses two distinct activation pathways, designated the classical and the alternative pathways (V. M. Holers, In Clinical Immunology: Principles and Practice, ed. R. R. Rich, Mosby Press; 1996, 363-391). The classical pathway is a calcium/magnesium-dependent cascade which is normally activated by the formation of antigen-antibody complexes. The alternative pathway is a magnesium-dependent cascade which is activated by deposition and activation of C3 on certain susceptible surfaces (e.g. cell wall polysaccharides of yeast and bacteria, and certain biopolymer materials). Activation of the complement pathway generates biologically active fragments of complement proteins, e.g. C3a, C4a and C5a anaphylatoxins and C5b-9 membrane attack complexes (MAC), which mediate inflammatory activities involving leukocyte chemotaxis, activation of macrophages, neutrophils, platelets, mast cells and endothelial cells, vascular permeability, cytolysis, and tissue injury.\nFactor D is a highly specific serine protease essential for activation of the alternative complement pathway. It cleaves factor B bound to C3b, generating the C3b/Bb enzyme which is the active component of the alternative pathway C3/C5 convertases. Factor D may be a suitable target for inhibition, since its plasma concentration in humans is very low (1.8 μg/ml), and it has been shown to be the limiting enzyme for activation of the alternative complement pathway (P. H. Lesavre and H. J. Müller-Eberhard. J. Exp. Med., 1978; 148: 1498-1510; J. E. Volanakis et al., New Eng. J. Med., 1985; 312: 395-401).\nThe down-regulation of complement activation has been demonstrated to be effective in treating several disease indications in animal models and in ex vivo studies, e.g. systemic lupus erythematosus and glomerulonephritis (Y. Wang et al., Proc. Natl. Acad. Sci.; 1996, 93: 8563-8568), rheumatoid arthritis (Y. Wang et al., Proc. Natl. Acad. Sci., 1995; 92: 8955-8959), cardiopulmonary bypass and hemodialysis (C. S. Rinder, J. Clin. Invest., 1995; 96: 1564-1572), hyperacute rejection in organ transplantation (T. J. Kroshus et al., Transplantation, 1995; 60: 1194-1202), myocardial infarction (J. W. Homeister et al., J. Immunol., 1993; 150: 1055-1064; H. F. Weisman et al., Science, 1990; 249: 146-151), reperfusion injury (E. A. Amsterdam et al., Am. J. Physiol., 1995; 268: H448-H457), and adult respiratory distress syndrome (R. Rabinovici et al., J. Immunol., 1992; 149: 1744-1750). In addition, other inflammatory conditions and autoimmune/immune complex diseases are also closely associated with complement activation (V. M. Holers, ibid., B. P. Morgan. Eur. J. Clin. Invest., 1994:24:219-228), including thermal injury, severe asthma, anaphylactic shock, bowel inflammation, urticaria, angioedema, vasculitis, multiple sclerosis, myasthenia gravis, membranoproliferative glomerulonephritis, and Sjögren's syndrome.\nThere is a need for antibody therapeutics in the field of complement-mediated disorders, and humanized anti-Factor D antibodies, and antibody variants thereof, and fragments thereof (e.g. antigen-binding fragments), of the present invention provide high affinity antibodies useful to meet this need."}
{"text": "1. Field of the Invention\nThis invention relates to a wraparound splice case to insulate and protect a splice in a cable.\n2. Discussion of the Prior Art\nU.S. Pat. No. 4,142,592 to Brusselmans discloses an assembly for protecting and insulating a junction such as a splice between cables. The assembly includes a relatively rigid longitudinally split tube which is positioned about the junction. The assembly also includes relatively rigid support means, such as rings, which are positioned at the end of the tube. A plurality of flexible prongs, which may be attached to or formed integrally with said support means, project longitudinally at spaced intervals around the circumference of at least one end of the tube. A heat shrinkable sleeve of length greater than the length of the tube, support means and prongs, is disposed about the tube, support means and prongs. At least a part of the inner wall of the heat shrinkable sleeve is provided with a sealant so that on shrinking, the end regions of the sleeve overlying the free ends of the prongs are shrunk down over the conduits (or cables) and are sealed thereto. The heat shrinkable sleeve can be of any suitable material.\nWraparound splice case assemblies, identified as XAGA 800 and XAGA 900, are available commercially from Raychem Corporation. The wraparound assembly includes an inner tubular member, which can be two metal half-cylinders, and a heat shrinkable sleeve of a length greater than the length of the inner member. The wraparound assembly is positioned about the cable splice. The heat shrinkable sleeve is fastened with a closure rail device, such as that described in U.S. Pat. No. 3,455,336 to Ellis and is then heat shrunk over the inner member and underlying cable. The assembly is provided in kit form with the metal half-cylinders attached to the inner surface of the heat shrinkable sleeve with two strips of pressure sensitive adhesive strips.\nU.S. Pat. No. 3,875,325 to C. V. Anderson et al. discloses a telephone splice closure. The closure consists of two mating shells which have tapered mouth portions to accommodate a range of cable outside diameters. The two mating half-shells are hinged together to make a one piece cover. The non-hinged side of the half shells are provided with ridges which are fastened together with a channel after the closure is positioned about cable."}
{"text": "Gate valves are located intermittently along the length of gas and oil pipelines. Such valves are adapted to generally remain open but under dislocations in the fluid flow within the pipeline caused by, for example, a leak in the pipeline, the valves are each adapted to close thus shutting off the flow of oil or gas until the dislocation is located and repaired.\nGate valves and associated actuator systems to perform such functions are known. Such an actuator system and gate valve is disclosed in our U.S. Pat. No. 4,423,758 to Ellett entitled EMERGENCY SHUT DOWN DEVICE. Certain new and inventive improvements have been made, however, in the gate valve actuator system there disclosed.\nIn the apparatus disclosed in the aforementioned reference, the preload in the compression spring within the housing of the actuator is constant. In order to change this preload in the spring, which may be required, for example, when the actuator is used with a different valve, it was necessary to disassemble the Ellett device to vary the spring loading. The disassembly was difficult and time consuming.\nA further disadvantage with the Ellett apparatus related to the position of the piston after the preload in the compression spring was changed. When a spacer was added between the pull tube and the end plate, the piston would move upwardly with the pull tube with the result that the stroke necessary to completely close the valve would be shortened by a like amount and this reduced length stroke might not fully close the valve. To compensate for this stroke reduction, a spacer was added between the piston and the pull tube. Because the stop tube is integral with the pull tube, however, and does not move relative to the pull tube, it was necessary to either provide a new pull tube and stop tube or to reduce the length of the stop tube by cutting or machining. This was also costly and unnecessarily time consuming.\nYet a further disadvantage in the aforementioned Ellett apparatus related to the necessity in actuators to have different stroke lengths for different valves. In the Ellett apparatus, there was no adjustment for the stroke possible and, therefore, there was little universality of the actuator for valves of various sizes.\nYet a further undesirable feature of existing actuators relates to the use of the indicator rod which indicates whether or not the valve is open or closed. It is desirable that the indicator rod attached to the piston be of a configuration which allows the valve to be mechanically maintained in an open position in the absence of pressure in the cylinder.\nYet a further disadvantage with existing actuators relates to removing the actuator from the valve while maintaining the valve in the open position there being a large force exerted on the piston by the compression spring. Disassembly of the actuator and, in particular, removing the hub from the end plate under such conditions can be dangerous."}
{"text": "1. Field of the Invention\nThe present invention relates to machinery for picking up and depositing raw dough onto pans, and more particularly to a machine which functions automatically and continuously to deposit rectangular strips of raw dough onto pans to form a dough product which can be easily baked as a sandwich-type bun.\n2. Description of the Prior Art\nThere are a number of prior art arrangements for automatically cutting and depositing gobs of dough into pans, such that the dough may subsequently be baked. For example, Roth et al disclose, in U.S. Pat. No. 3,596,613, automatic machinery for continuously cutting and depositing gobs of raw dough into pans.\nIn U.S. Pat. No. 2,018,950, Couture discloses an assembly line technique for making sandwich-type biscuits. Other prior art references of interest include the following U.S. Pat. Nos. 3,397,655 to Valadez; 2,970,268 to Doolin; 2,780,182 to Rand; 2,602,392 to Pankin; 3,111,913 to Maadek; 2,352,447 to Powers; 2,186,807 to Newberger; 2,014,384 to Kruttschnitt; 2,570,374 to Pompa; and U.S. Pat. No. Re. 22,399 to Fonken."}
{"text": "The present invention relates to a method of heat treatment of a low melting-point metal for production of spheroidal bodies formed of a low melting-point metal or alloy (hereafter collectively referred to as the low melting-point metal) or heat treatment of conductor portions of electronic circuit boards used for electronic apparatus or land portions thereof where conductors or parts are mounted. More particularly, the present invention concerns a method of treating electrodes of chip parts, such as chip resistors, i.e., a kind of electronic parts, which is capable of contributing to the production of lightweight, thin, and compct electronic equipment.\nHitherto, a heat treatment method of this type has been conducted as described below. FIG. 1 is a diagram schematically illustrating an arrangement of a conventional jig for forming low melting-point metal balls. In the drawing, reference numeral 1 denotes a graphite plate; 2, a semispherical hole provided in the graphite plate 1; and 3, a low melting-point metal pellet placed in the hole 2. To describe a method of forming low melting-point metal balls using the jig shown in FIG. 2, the graphite plate 1 having a multiplicity of the semispherical holes 2 is prepared, the low melting-point metal pellets 3 are placed in the holes 2, and the graphite plate 1 thus placed in the graphite plate 1 is immersed in a high boiling-point liquid (not shown) having a temperature higher than the melting point of the low melting-point metal. When the low melting-point metal pellets 3 have melted, and lower portions thereof have become semispherical in conformity with the semispherical holes 2 provided in the graphite plate 1 and upper portions thereof have become semispherical due to the surface tension, the graphite plate 1 is removed from the high boiling-point liquid of a high temperature and the metal balls are allowed to solidify so as to obtain the low melting-point metal balls.\nMeanwhile, as for a circuit board used in an electronic apparatus, one having a structure as shown in FIG. 2 is used. FIG. 2 is a perspective view of a circuit board used for, for instance, a composite part. In the drawing, reference numeral 4 denotes an insulating substrate; 5, a land to which an electronic part such as a chip is soldered; 6, a conductor for connecting lands; and 7, a land for soldering a terminal thereto. These lands 5, 7 and conductors 6 are disposed on the substrate 4. In a circuit board such as the one shown in FIG. 2 which is currently used, copper foils are in most cases used as the lands 5, 7 and conductors 6. Recently, a circuit board in which a film of a low melting-point metal, such as solder, is coated on the aforementioned copper foils has come to be used to improve the reliability of soldering. As a method of coating with this low melting-point metallic film, an electroplating method, a chemical plating method, or a method of immersing in a low melting-point metal is generally adopted. A circuit board coated with the low melting-point metallic film by such a method is used as it is without being subjected to heat treatment.\nIn such a conventional heat treatment method, when low melting-point metal balls shown in FIG. 1 are produced, the low melting-point metal pellets must be arranged with adquate intervals provided therebetween in such a manner that the metal pellets will not come into contact with each other when the metal pellets are melted. In addition, a pallet, such as the graphite plate having semispherical holes, is required to ensure that the low melting-point metal pellets will not adhere thereto, and the pellets must be placed into the holes one by one. This constitutes a major factor impeding mass production. Furthermore, although the low melting-point metal pellets are placed in the semispherical holes and are melted to produce the low melting-point metal balls, it is difficult to have the radius of curvature of a spherical surface formed by the semispherical hole conform with that of a spherical surface formed in an upper portion of the hole by the surface tension, so that there is a drawback in that balls with a desirable spherical shape cannot be obtained. Moreover, in terms of facilities, both an equipment for placing the low melting-point pellets into the semispherical holes in the pallet and another equipment for removal after heat treatment are required. As for heat treatment facilities, since the low melting-point pellets must be subjected to heat treatment together with the pallet, a large heating capacity is required, so that large-scale facilities are required as a whole.\nIn addition, a method of producing spheroidal metallic particles is disclosed in Japanese Patent Unexamined Publication No. 56-84401. This method of producing spheroidal metallic particles is effected as follows: As shown in FIG. 3, a heat insulating glass tube 8 is filled with a liquid 10, such as glycerin, which is capable of maintaining a liquid state even at a temperature higher than the melting point of metallic particles 9 and has reducing properties. Heaters 11, 11a, and 11b are wound around an outer periphery of the glass tube 8, an uppermost layer of the liquid 10 in the glass tube 8 is maintained at a temperature below the melting point of the metallic particles 9 by means of the uppermost heater 11. In addition, an intermediate layer of the liquid 10 is held at a temperature above the melting point of the metallic particles 9 by means of the intermediate heater 11a. Furthermore, the temperature of the liquid 10 is set by the lowermost heater 11b such as to gradually decline from the lower layer of the intermediate layer to the lower end of the glass tube 8 below the melting point of the metallic particles 9. Subsequently, the metallic particles 9 are poured into the liquid 10 from a hopper 12 above the glass tube 8.\nNamely, in the uppermost layer whose temperature is held below the melting point of the metallic particles 9, the metallic particles 9 poured into the liquid 10 are separated from each other as each of them becomes completely wet with the liquid 10. The separated metallic particles 9 then fall further in the liquid 10, are heated and melted in the liquid 10 in the intermediate layer whose temperature is held above the melting temperature of the metallic particles 9, and are formed into spheroidal metallic particles 9a by virtue of the surface tension. The molten metallic particles 9a thus formed into the spheroidal shape further fall in the liquid 10, and solidify as they are in the spheroidal shape in the lowermost layer portion of the liquid 10 whose temperature is set below the melting point of the metallic particles 9 (9a). The metallic particles 9a thus solidified are collected in a collector 13 provided at a lower end of the glass tube 8. Subsequently, these spherical metallic particles 9a are taken out, and glycerin or the like adhering to their particle surfaces is washed away to obtain granular metallic particles having smooth, glossy surfaces which are free from uneveness and cracks.\nAccording to this method, however, there have been the following problems. First, the arrangement is very complicated since the temperature gradient of the liquid is such that, as compared with the melting point of the metallic particles, the liquid temperature gradient provided from the upper portion of the glass tube downward is in twin set to become a temperature below the melting point, another temperature above the melting point, and still another temperature below the melting point (and another temperature further below that temperature). Furthermore, the heating devices become complicated in order to form such a temperature gradient, a number of such units are required. In other words, this method has had a serious drawback in that the temperature gradient from the uppermost portion to the lower portion of the glass tube is complicated, and means for controlling this temperature gradient are made complicated. Thus, there has been a problem in that efficiency in mass production is extremely poor.\nSecondly, when a copper foil of a circuit board, such as the one shown in FIG. 2, is subjected to a low melting-point metallic film treatment, in the method of coating the copper foil with the low melting-point metallic film by electroplating or chemical plating, the wettability of solder is good and soldering can be effected without any trouble insofar as soldering is effected within a fixed period of time after plating. According to these plating methods, however, the surfaces of the low melting-point metallic films are coarse, and their surface areas are very large. For this reason, these films are liable to adsorb foreign substances and gases, and if they are stored for a long period of time, the low melting-point metal surfaces undergo chemical change such as oxidation. Hence, such a method disadvantageously has a large possibility of causing faulty soldering at the time of soldering electronic parts or the like. Meanwhile, according to a method of immersing in a molten low-temperature metal, the thickness of the low melting-point metallic film becomes quite nonuniform. Hence, this metod has a serious drawback in that faulty mounting of parts occurs due to the excessive thickness when parts such as chips are actually mounted on the circuit boards.\nIncidentally, if the low melting-point metal plated film is constituted by a glossy plating, solderability is poor since impurities (organic substances) are contained therein, so that this method is also fraught with a serious drawback."}
{"text": "1. Field of the Invention\nThe present invention relates to a mobile communication system, and more particularly, to a method of transmitting or receiving a data packet in packet data communication system using a hybrid automatic repeat request.\n2. Discussion of the Related Art\nGenerally, a hybrid automatic repeat request (HARQ) is used for transmission of packet data in a mobile communication system.\nThat is, the object of the HARQ is to improve the reliability and throughput during data transmission in the communication system by combining an automatic repeat request (ARQ) with a forward error correction (FEC).\nThe object of the ARQ is to improve the reliability in a manner that a receiving end requests repeat of the same information until the initially sent information is received in the receiving end without error. The object of the FEC is for the receiving end to compensate for the error occurring due to the channel environment using an error correction code.\nIf the channel environment is good and the error occurring frequency of the given information is low, only the use of the ARQ will suffice for the reliability.\nHowever, if the channel environment gets worse, the error occurring frequency of the given information becomes greater, and thus the frequency of repeat request also becomes greater.\nThis deteriorates the throughput of the system. Thus, it has been proposed to use both the ARQ and the FEC, which is called the HARQ.\nAs a kind of the HARQ, there is a method that uses an incremental redundancy.\nAccording to this method, a transmitting end first transmits information encoded with a high coding rate to a receiving end, encodes the information with the coding rate lowered whenever the repeat request is received from the receiving end, and then transmits only the additional redundancy bits to the receiving end. The receiving end combines the received redundancy bits with the previously sent information to decode the information.\nThese additional redundancy bits are bits added for error correction or detection of the previously transmitted data packet.\nTo do so, the receiving end can obtain the gain according to the combining. Also, since the combining coding rate of the packet data is gradually reduced whenever the repeat is performed, the transmitting end can adaptively give the redundancy to the data packet according to the change of the channel environment.\nCurrently, in a high data rate (HDR) mobile communication system, the incremental redundancy of the HARQ, particularly a synchronous incremental redundancy (SIR) is used.\nAccording to the SIR method, as shown in FIG. 1A, information to be transmitted is encoded and repeated to construct one packet, and then the data packet is divided into sub packets of a fixed size for transmission.\nSpecifically, the transmitting end transmits a sub packet among sub packets that constitute a data packet. If the receiving end performs decoding of the transmitted sub packet without error and transmits an acknowledgment (ACK) to the transmitting end, the transmitting end transmits no further sub packet.\nIf the receiving end fails to perform the decoding of the transmitted sub packet and transmits a non-acknowledgment (NACK) to the transmitting end, the transmitting end transmits the next sub packet.\nMeanwhile, in order to heighten the transmission rate in the changing channel environment, it is advantageous to change the length of the respective sub packets according to the changing channel environment.\nAccordingly, in order to implement this, a method that uses a variable sub packet length has been proposed as shown in FIG. 2A.\nDuring the process of dividing a data packet into several sub packets and transmitting the sub packets, this method enables the transmission of the respective sub packets to be performed with different sub packet lengths according to the channel environment.\nAt this time, there is provided a control channel for notifying information on the length of the sub packet, sub packet identifier (ID), etc.\nIf the transmitting end transmits the packet data to the receiving end through the above-described process, the receiving end extracts information on the sub packet through the control channel to receive the sub packet.\nThen, the receiving end performs the decoding of the data packet by concatenating/combining the received sub packet with the data packet previously sent from the transmitting end using the extracted information.\nIn case of using the sub packets of the fixed length as shown in FIG. 1A, the length of all the sub packets is fixed, and thus the receiving end can grasp the transmission start point of the currently received sub packet if it knows the order of the sub packets even when it cannot receive any sub packet. Thus, it can perform the decoding by regenerating the packet with other sub packets (or even with only one sub packet) as shown in FIG. 1B.\nThus, if there is information on the order of the sub packet currently transmitted, the receiving end can recognize the transmission start point of the sub packet. In case of using the SIR method in the HDR system, the receiving end can always recognize information on when its own sub packet is received after the first sub packet is received, and thus does not require other control information in the event that the received sub packet is not the fifth sub packet. However, in case of using an asynchronous incremental redundancy, the transmission time of the next sub packet may randomly change even if the transmission time of the first sub packet is determined, and thus all the sub packets should be transmitted to the receiving end along with the information on the sub packet ID and the sub packet length that is the control information on the sub packet. In case that the length of the respective sub packet is fixed, the receiving end can recognize the transmission start point of the respective sub packet even if the transmitting end just transmits the sub packet IDs in order.\nHowever, according to the method that uses the sub packets of the variable length as shown in FIG. 2A, the length of the sub packet is changed whenever the sub packet is transmitted, and thus if the control channel that includes the sub packet information is not received, the start point of the data packet cannot be recognized only through the respective sub packet ID.\nThat is, it cannot be recognized where the currently received sub packet starts. For this reason the variable length type method cannot decode the data packet even if only one sub packet is missing.\nSpecifically, if a sub packet 2 is not received and a sub packet 3 is received, there is no information on the previous sub packet 2 during the process of concatenating the sub packet 1 and the sub packet 3, and thus it cannot be recognized of which type between two types as shown in FIGS. 2B and 2C the sub packet 3 is received.\nIn the sequential transmitting method wherein the transmission of the next sub packet starts at a position where the transmission of the previous sub packet ends, if the receiving end misses even one sub packet, it cannot recognize the transmission position of the incoming sub packet after the missing sub packet, and thus the decoding cannot be performed.\nTo solve the above-described problem, a method is used as follows. That is, in transmitting variable sub packets, each sub packet is transmitted from the fixed start position in the encoded packet to the receiving end, and the information on the length and the starting position of the sub packet is transmitted to the receiving end through the control channel. The receiving end extracts the information on the sub packet through the control channel, and performs the decoding by concatenating/combining the received sub packet with the previously received sub packet using the extracted information.\nHowever, in case of using this method, a certain portion in the encoded packet may be repeatedly transmitted, while another portion may not be transmitted, resulting in that the coding gain according to the HARQ cannot be obtained sufficiently."}
{"text": "This invention relates generally to foldable garment support frames and, more particularly, to a new and improved foldable garment support frame which firmly supports and maintains garments in a lightly stretched condition to minimize wrinkling that otherwise results from folding and packing such garments in luggage.\nFolding garment fixtures are known which are adapted to aid the traveler in folding and maintaining garments in a neat and relatively wrinkle-free condition. Such garment frames are most useful for packing jackets, trousers and related clothing accessories into the limited confines of a suitcase. These garment frames are intended to inhibit a substantial amount of the wrinkling which otherwise would form in garments that have been folded tightly in a suitcase for long hours.\nOne known variety of foldable garment fixture comprises a conventional coat-and-trouser hanger suspended from an anchor point on a loop frame formed of metal rod, the frame including a swingable folding bar. To use such a fixture, the garment, typically a jacket and trousers, is arranged on the hanger and suspended from the anchor point. The folding bar is swung into place, laterally disposed across the garment about midway along its length. The garment is then folded around the folding bar back onto itself, being careful that the garment, and particularly any jacket sleeves, are neatly arranged to avoid wrinkling. Such a fixture is usually designed for use with a relatively large suitcase.\nIn a variation of the foregoing garment fixture, the garment may be folded about a second laterally disposed folding bar, with the result that the garment is folded approximately in thirds for packing in smaller suitcases, sometimes denoted \"overnight\" cases.\nIn using either of these foldable garment fixtures, it will be noted that care must be taken to maintain the garments in neat order when folding the assembly. Further, it is apparent that firm support is not provided over the major portion of the garment area which is spanned between the hanger and the folding bar or bars. In other words, the garment is still susceptible to being crushed or wrinkled by other objects which are likely to be placed next to or on top of the garment in packing because the garment is not backed by a firm, flat structure. One solution to this problem is to not pack such objects in the same luggage compartment with the subject garment, but this can result in a substantial waste of space.\nAdditionally, with these known foldable garment fixtures, there are no special means provided for maintaining the garment relatively tautly wrapped about the folding bars. If the garment should slip or loosen on the fixture, it is more likely to be wrinkled. On the other hand, a garment which is maintained in a lightly stretched condition is less subject to wrinkling.\nAccordingly, a need is recognized for a foldable garment support frame which can be used easily for folding and packing garments, will firmly support the entire area of the garment and is effective to maintain the garment in a lightly stretched condition on the frame. The present invention clearly fulfills this need."}
{"text": "Memory systems using \"hard disks\" to store information in annular tracks concentrically arranged in magnetic material formed on a disk surface are well known in the prior art. The information (often called \"data\") is stored in magnetic domains and each domain typically stores an encoded data bit. A major advantage of this type of memory system is that information can be easily stored on a selected track on a disk surface as the disk rotates at high speed and can be randomly accessed at high speed to store information or to read out information stored on a given track. A read-write head or transducer having a small coil and a magnetic core is moved across the disk surface to the track on which the desired information is to be stored, or from which the desired information is to be retrieved while the disk is spinning. The read-write head stores or detects data in a selected data track on the disk surface. Ideally, the head rides on a thin layer of air generated by the spinning disk. It is then said to fly over the disk surface. However, due to imperfections in the disk surface or to other operating factors such as system vibrations, the head may come into undesired contact with the magnetic film during flight while the disk is spinning, causing instant friction between the head and the magnetic film which may ultimately lead to the destruction of the film and to a head crash. Also, during start/stop there is intimate contact between the head and the disk surface with associated friction which may lead to excessive wear, head crash, or stiction. Stiction may be defined as head/disk adhesion strong enough to require manual removal of the head, resulting in severe marks on the disk surface. Stiction is often encountered when using relatively thick liquid lubricant or when the memory system is operating under conditions of high relative humidity.\nIn the prior art various techniques have been used to reduce the friction between the head and the magnetic film and the associated damage caused by it.\nIn memory systems which utilize \"flexible\" media such as magnetic tape or floppy disks, the space between the moving medium and the read/write head is much smaller than in hard disk systems, and there is continuous intimate contact between the medium and the read/write head. In both hard and flexible media, a magnetic film is formed on a substrate. For a hard disk, the substrate is typically an aluminum alloy with a relatively thick and hard undercoat formed thereon. The magnetic film is formed on the undercoat and is typically a cobalt phosphorus alloy or particulate oxide. The technique used to reduce friction or contact damage depends on the composition and quality of the magnetic film.\nOne method of forming a magnetic film on a disk substrate involves suspending iron-oxide power in a liquid media containing resins and coating the disk surface with the liquid. The disk is then cured in a strong magnetic field to provide the desired anisotrophy of magnetization. Such particulate magnetic films are typically formed to a thickness of between 0.7 and 5.0 microns, although recently particulate films having a thickness approaching 0.5 micron have been employed. These particulate iron-oxide magnetic films have a porous surface, and for these surfaces, one antifriction coating technique consists in applying a thin film of an organic lubricant, i.e., an organic liquid or wax such as Fomblin or Bareco. These organic lubricants are not very satisfactory because they tend to spin off the disk when the disk is rotated at operating speeds, e.g., 3600 rpm.\nThe trend in magnetic films has been to thinner films, because generally speaking, the thinner the magnetic film, the greater the information storage density, and to metal alloy film, as opposed to particulate oxide films. Particulate oxide films are porous and have low remanance due to their relatively small volume fraction of magnetic material in the resin. Since the magnitude of the read back signal is directly related to the remanance of the film, particlate oxide films are necessarily thick relative to alloy films of the same remanence.\nThe quest for thin magnetic films has led to several deposition techniques. Thin films typically comprise magnetic alloys of cobalt with such elements as phosphorus, nickel, chromium, or tungsten. Methods of deposition of thin magnetic alloy films include electroless, electroplating, sputtering, evaporation and other vapor deposition techniques. Thin magnetic films may also be formed by vapor deposition or sputtering or a magnetic metal oxide such as iron oxide in the form of Fe.sub.3 O.sub.4 or Fe.sub.2 O.sub.3 or by oxidation of metallic iron or by thermal decomposition of iron salts. Thin magnetic films typically have a thickness of 600-2000 .ANG., although even thinner films are extant, some having a thickness in the range of 300 .ANG.. These thin magnetic films generally have a smooth, nonporous surface, and the prior art liquid or waxy lubricants that were marginally satisfactory for particulate films, are less satisfactory for smooth, thin films since a liquid or waxy lubricant spins off the disk after a short period of operation.\nPowdered non-metallic solid lubricants, such as graphite and molybdenum disulfide have been applied to plated thin magnetic films by rubbing under pressure, but have proven to be unsatisfactory. (See Yanagisawa, U.S. Pat. No. 4,398,562, issued June 28, 1983). In general, these lubricants are made of loose particles which often detach and break, leading to a head crash. The particles are too large, e.g. 3 microns, average diameter, for an adequate thin layer of lubricant, and they present poor adhesion to metallic magnetic films.\nSince no satisfactory antifriction lubricants have been devised for thin magnetic films, the trend has been to use an extremely hard protective coating, such as plated rhodium and nickel-tin, which have a hardness in the range of 500-1000 Kg/mm.sup.2, and sputtered carbon, which has a hardness in the range of 3,000-5,000 Kg/mm.sup.2, as shown by R. J. Gambino and J. A. Thompson. Solid State Communications, Vol. 34 pp. 15-18 (1980). While these hard coatings tend to protect the magnetic medium from the head riding above the disk they have several drawbacks: one, they are generally expensive to implement; two, since they are so extremely hard, they tend to cause wear of the head riding on their surface. This tends merely to shift the wear problem from the disk to the head (and increasing the hardness of the head tends to shift the problem back to the disk).\nIn the case of sputtered carbon, the sputter deposition equipment is every expensive, the throughput low, and it is difficult to control its adhesion to the magnetic film and its crystalline character.\nRhodium, which is often taught, is still unsatisfactory in its performance. To assure effective protection, a thick layer of rhodium, typically on the order of about ten micro inches, has been employed. See U.S. Pat. No. 3,767,369, \"Duplex Metallic Overcoating\", issued Oct. 23, 1973, which is incorporated herein by reference. The U.S. Pat. No. 3,767,369 proposes an alternative to rhodium comprising a duplex of an at least 2 micro inches thick layer of an alloy of nickel and tin overlaid by a layer of rhodium at least 2 micro inches thick. The nickel-tin layer contains a proportion of nickel of about 35% by weight. Again, this nickel-tin alloy undercoat has a Vickers hardness of 500 to 900. More recently, a layer of rhodium 750 .ANG. thick has been suggested as a protective layer in combination with a thin layer of liquid lubricant, Dupont 804. See Rossi, et al, \"Vacuum Deposited Thin Metal Film Disk\", J. Appl. Phys., 55 (No. 6) 2254 (1984), which is incorporated herein by reference. Such a combination ignores the long-term reliability problem due to the spin-off of the lubricant. Also, the total thickness of the overcoat, including chromium, rhodium, and liquid lubricant is about 1,000 .ANG.."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display (LCD) device, and more particularly, to an LCD device that can reduce a fabrication cost and can prevent a decrease in picture quality caused by unfilled pixel cells.\n2. Discussion of the Related Art\nA liquid crystal display (LCD) device controls an optical transmission ratio of liquid crystal molecules using an electric field to thereby display an image. In general, an LCD device is provided with a liquid crystal panel having liquid crystals between two glass substrates, a matrix of liquid crystal cells, with switching devices for respectively changing signals in the liquid crystal cells, a driving circuit for driving the liquid crystal panel, and a backlight unit for directing light onto the liquid crystal panel.\nRecently there have been efforts to reduce the number of signal lines or the number of circuit components of the liquid crystal panel to develop a thin and light, as well as low cost liquid crystal display device. For example, Korean patent publication No. 10-2003-0039972 discloses a flexible printed circuit and a liquid crystal display device without a gate printed circuit board that can reduce the panel size and the production cost of the liquid crystal display device.\nHowever, because a plurality of data and gate driver integrated circuits, data printed circuit boards, and the like are provided, the flexible printed circuit and the liquid crystal display device even without having a gate printed circuit board have problems in that the reduction of the size and the production cost is limited, and power consumption is high.\nFor the related art LCD device, if the ambient temperature of liquid crystal panel is low (for example, below −5° C.), the charging property of pixel cell is lowered so that it causes the deterioration of picture quality due to unfilled pixel cells."}
{"text": "In conventional technologies, the structure of a TFT-LCD (Thin Film Transistor Liquid Crystal Display) mainly comprises two glass substrates and a liquid crystal layer disposed therebetween, wherein a surface of the upper glass substrate has a color filter mounted thereon, the color filter includes color filter units and black matrix; and a surface of the lower glass substrate has thin-film transistors and pixel electrodes mounted thereon, therefore it is generally called TFT Array Substrate.\nIn a thin-film-transistor array substrate, the structure of a sub-pixel comprises a pixel electrode, a thin-film transistor and a storage capacitor, wherein a gate electrode of the thin-film transistor is connected a scan-signal line, a source electrode thereof is connected to a data-signal line, a drain electrode thereof is connected to the pixel electrode; and the storage capacitor is connected to the pixel electrode. By means of the scan-signal line, the gate electrode of the thin-film transistor is forced with a voltage, which is capable of conducting the thin-film transistor. The source electrode of the thin-film transistor then, by means of the data-signal line, receives a data signal and transmits the data signal to the drain electrode, and further writes it in the pixel electrodes and stores it in the storage capacitor.\nBesides the storage capacitor, many parasitic capacitances further exist in the actual structure of the thin-film transistor further, those parasitic capacitances are unnecessary to the thin-film transistor and may cause electronic losing and affect the operation characteristic of the thin-film transistor.\nWith reference to FIG. 1, FIG. 1 is a scheme diagram of a partial structure of a conventional thin-film-transistor array substrate. As shown in FIG. 1, a thin-film transistor 9 is disposed close to a point of the junction of a scan-signal line 8 and a data-signal line 7, wherein, the thin-film transistor 9 includes a gate electrode 90, a semiconductor layer 91, a drain electrode 92 and a source electrode 93, wherein, the gate electrode 90 is a portion of the scan-signal line 8; the semiconductor layer 91 is mounted upon the gate electrode 90 with an isolation layer placed therebetween (not shown in the figure); the drain electrode 92 is disposed on the semiconductor layer 91; the source electrode 93 extends from a side of the data-signal line 7 along the edge of the semiconductor layer 91 and then surrounds the drain electrode 92. Parasitic capacitances exist between the drain electrode 92 and the gate electrode 90.\nWith further reference to FIG. 2, in order to decrease the parasitic capacitances that exist between the drain electrode 92 and the gate electrode 90, the gate electrode 90 further has an opening 900 corresponding to the drain electrode 92, so as to reduce an overlapping area of the gate electrode 90 to the drain electrode 92, and thereby decreases the parasitic capacitances formed therebetween.\nHowever, the opening 900 of the gate electrode 90 will reduce the area of the scan-signal line 8, and thereby increase the numeral value of the resistance of the scan-signal line 8. The transmission speed of the scan-signal line 8 will be affected due to the increased resistance.\nTherefore, it is necessary to provide a thin-film-transistor array substrate and a manufacturing method thereof to overcome the problems existing in the conventional technology."}
{"text": "Optical fiber grating devices having outer coatings of variable thickness metal are highly useful in optical communication systems. In these devices variable thickness thin metal films are used as thin film resistive heaters or strain-relief coatings for altering in controlled ways, the optical properties of the grating. For example, such devices can provide Bragg gratings with adjustable chirp for use as tunable dispersion compensating elements. See U.S. patent application Ser. No. 09/183,048 filed by Eggleton el al. on Oct. 30, 1998 and entitled \"Optical Grating Devices With Adjustable Chirp.\"\nThe ability to form on-fiber metal coatings with complex but precisely defined variations in thickness is critical to producing these and other types of active grating devices. One method for forming such coatings involves electrodeposition of metal onto a fiber while it is pulled from a planting bath at a controlled rate. Although this technique can form tapered coatings with a range of thicknesses, it has the disadvantages that (i) non-monotonic variations in thickness are not straightfonvard to produce, (ii) a seed layer, typically deposited by electron beam evaporation or electroless planting, is required to initiate electrodeposition, and (iii) formation of ultrathin (&lt;1 .mu.m) coatings with precise variations in thickness is challenging. Accordingly there is a, need for an improved method for providing optical fiber devices with metal coatings of controlled variable thickness."}
{"text": "1. Field of the Invention\nThe present invention relates to a dietary supplement and method for use as a probiotic(immuno-stimulant) and for alleviating the symptoms associated with irritable bowel syndrome.\n2. Description of Related Art\nThe use of dietary supplements for alleviating specific symptoms associated with particular human health problems is well known. Going back to ancient times, references are made to various dietary foods, herbs, and other naturally produced substances that are associated with improving the human health condition.\nMany times, a person has poor health conditions based on dietary deficiencies, such as vitamins, minerals, and other natural elements that are essential for good human health. Many times, the dietary supplements available can supply missing dietary minerals and vitamins that are essential for proper metabolism and function of the human body. Therefore, it is well known that improvements to human health and well-being can be achieved using dietary supplements, especially if they provide for deficiencies in a person's diet. The use of daily or weekly supplements supplying missing natural dietary elements such as vitamins and minerals can alleviate symptoms for improvement related to specific maladies of human health.\nThe human condition known as irritable bowel syndrome can best be defined as a chronic disorder of the bowels, resulting in variable abdominal discomfort or pain, constipation, diarrhea, cramps, nausea, belching, excess flatus, spasm, abdominal distention, tired and weak feeling, and mechanical irritation of the bowel. Medical treatment has traditionally entailed bland diet, mild sedatives, hydrophilic colloid laxatives, anti-cholinergics, and anti-diarrhea agents and are often associated with numerous side effects such as drowsiness, dry mouth, visual disturbances, and food intolerances.\nAloe is a tropical or subtropical plant with thick lance-shaped leaves with jagged edges and has been known for centuries for its medicinal and therapeutic properties useful in treating many different human conditions.\nThe aloe plant and extracts from aloe, aloe vera, and various specific aloe extracts from the natural plant can alleviate symptoms related to irritable bowel syndrome as shown in the following literature.\nTitle or Subject Authors Date Reference \"The Small Intestines Chikalo, I. & Sept. Quoted in \"The Function Affected by Bolovyeve, V. 1966 Aloes of Aloe Extract\" Ukraine Tropical Africa and Madagascar\" by Reynolds, G. W. \"Aloe Vera Gel in Blitz, J., 1963 J. Amer. Peptic Ulcer Therapy; Smith, J., & Osteopathic Preliminary Report\" Gerard, J. Assoc. 62 731- 735 Aloenin and Aloe-Ulcin Yamamoto, I. 1970 J. Med. Soc. From A arborescens Toho Jpn. 17 Inhibit Gastric 361 Secretion Aloe-Ulcin Equals Mg Hirata, T. & 1977 Naturforsch 32c Lactate Suga, T. 731 Aloenin and Aloe-Ulcin Hirata, T. & 1978 Bull. Chem. From A arborescens Suga, T. Soc. Jpan. 51 Inhibit Gastric 842 Secretion \"Effect of Orally Bland, J. 1985 Preventive Consumed Aloe Vera Medicine, Juice on March/April Gastrointestinal Function in Normal Humans\" \"Effects of Aloe Salto, H., 1989 Jpn. Yakugaku Extracts Aloctin A on Imanishi, K., Zasshi 109 (5) Gastric Secretion and Okabe, S. 335-339 on Experimental Gastric Lesions in Rats\" \"A Double-Blind Trial Odes, H. S. & 1991 Digestion -49 of a Celandin, Aloe Madar, Z. (2) 65-71 Vera and Psyllium Laxative Preparation in Adult Patients with Constipation\"\nFreeze-dried aloe and extract from the natural aloe plant have been determined to have a positive influence as a dietary supplement and probiotic to improve the overall health and condition of a person. Specifically, freeze dried aloe and aloe extracts have been found to have antibacterial and antiviral properties.\nThe following information is from UK Patent No. GB 2245143B, describing a freeze-dried aloe process that could be used to prepare aloe for use in the present invention. Aloe has been extensively used as medicine. In U.S. Pat. No. 5,455,033 to Silverman, the soothing action of aloe vera extract was disclosed and utilized in an topical anti-inflammatory composition. In addition, the Silverman patent reported inulin, an ingredient in Applicant's present invention, as being an activator in the immune system and the root extract demonstrating anti-viral activity against influenza, herpes and vesicular stomatis viruses. Reference to aloe powders is also made in the Japanese Pharmacopeia. Certain low, medium, and high molecular weight compounds have been separated from aloe. For example, it is known that aloe arborescens var. natalensis (referred to hereinafter as \"Kidacki aloe\") contains polysaccharides such as aloetin, alcenin, aloeursin or D-mannose, aloemannan or aloe albonacite; that aloe vera contains polysaccharides; and that Cape aloe contains anthraquinonic ingredients such as aloin or aloe-emodin. Bioactive factors were separated from freeze dried aloe in U.S. Pat. No. 5,902,796 to Shand. Mixtures of these active chemical substances identified, isolated, and stabilized from aloe vera leaves have been described in U.S. Pat. Nos. 4,735,935; 4,851,224; 4,917,890; 4,957,907; 4,959, 241; and 4,966,892. One group of these active chemical substances has been referred to as aloe vera mucilaginous polysaccharides, which are comprised of oligomers and polymers of carbohydrates. The pharmaceutical application of these mucilaginous polysaccharides and uses of aloe products have been described in U.S. Pat. Nos. 5,106,616; 5,118,673; 5,308,838; 5,441,943; and 5,443,830.\nIntensive attempts have been made to process aloe into medicines. Presently, as an officially allowed health food, food products prepared from Kidachi aloe have been extensively commercialized in Japan. Most of these products, however, are prepared by drying aloe leaves under the sun or by hot air and, at the time of tableting, up to about 10% of an additive such as Avicel.RTM. microcrystalline cellulose or cornstarch, is added as a binder to facilitate tableting.\nIt is noted that most ingredients isolated and purified from the aloe and found to be medically efficacious to animals by animal experiments, are polysaccharides, glycoproteins, or enzymes, all of which may exhibit activity in a stable condition through purification at lower temperatures. For this reason, it is generally thought to be desirable to apply as little heat as possible for the preparation of a health food or medicine from aloe.\nHowever, since the products now available on the market are prepared by drying aloe leaves under the sun or by circulation of hot air, the polysaccharides, glycoproteins, and enzymes in them will have been thermally degraded or oxidized in the course of drying.\nFurthermore, the additives for facilitating tableting, such as Avicel.RTM. or cornstarch, are added in an amount of several to tens of percents in the conventional aloe tablet products. Since a small amount of water is added under heating to knead the additives with the dried aloe mass, the resulting product undergoes further thermal degradation by such heating.\nSoluble fiber in the diet is well known for its salutary effects on gastrointestinal health and for being a very good fermentation substrate for intestinal flora. Such salutary effects include providing bulk to the stool, decreasing the pH of the gastrointestinal tract, producing volatile fatty acids, decreasing intestinal transit time, and beneficially influencing various blood parameters. Dietary fiber has also been shown to have a beneficial effect on cholesterol and lipid metabolism that results in decreased serum cholesterol, triglycerides, and phospholipids and an improved (increased) HDL to LKL ratio. A study on laboratory animals showed that adding fiber to the diet decreases the incidence of bacterial translocation, i.e. corrsing the intestinal barrier and entering systemic circulation. C. Palacio et al., Dietary Fiber: Physiologic Effects and Potential Applications to Enteral Nutrition, in Clinical Nutrition: Enteral and Tube Feeding (2d. ed., 1990). Nutritional and epidemiological studies have indicated that a general increase in the consumption of dietary fiber may play a role in preventing deleterious effects of oxygen free radicals that have been involved in such processes as aging, inflammation, and some disease processes. R. Kohen et al., Prevention of Oxidative Damage in the Rat Jejunal Mucosa by Pectin, 69 Br. J. Nutrition 789 (1993).\nInulin is a naturally occurring soluble fiber, a fructo-oligo saccharide composed of a mixture of oligomers of varying degrees of polymerization or molecular weights that occurs in numerous plants, including the dahlia tuber. It is not digested in the small intestine, but is fermented in the colon. The main effects of inulin on the digestive system are a decrease in the duration of the intestinal transit, a decrease in the level of glycemia, a decrease in the lipid content in the blood, a decrease in the pH in the colon, a decrease in the constipation phenomenon and a bifidogenic effect, for example. Thus, inulin can be fermented by bifidobacteria, which has the consequence of increasing the concentration of these bacteria at the level of the intestinal flora and of decreasing the concentration of enterobacteria, in particular Clostridiae, at the level of the intestinal flora.\nWhile prior art formulas as dietary supplements containing soluble fiber or aloe extracts are known and are generally suitable for their limited purposes, they possess certain inherent deficiencies that detract from their overall utility in restoring and maintaining gastrointestinal health. For example, a dietary supplement containing soluble dietary fiber without aloe extracts lacks antibacterial and antiviral activity. Similarly, a dietary supplement containing aloe without soluble dietary fiber lacks means for providing bulk to the stool, decreasing the pH of the gastrointestinal tract, producing volatile fatty acids, decreasing intestinal transit time, beneficially influencing various blood parameters, beneficially influencing cholesterol and lipid metabolism, decreasing the incidence of bacterial translocation, preventing deleterious effects of oxygen free radicals, and favoring the growth of beneficial bacteria in the gastrointestinal tract. Further, such prior art formulas fail to provide living intestinal bacteria that are beneficial for gastrointestinal health by providing an inhibitory effect on the growth of pathogenic bacteria, reducing the levels of toxic amines, and lowering the pH of the gastrointestinal tract.\nIn the prior art the production of inulin from plant materials such as Jerusalem artichoke, dahlia and chicory tubers is normally accomplished by using the following general procedure:\n1. Washing the tubers; PA1 2. Chopping, grinding or slicing the tubers; PA1 3. Extracting the inulin from the tubers with water; PA1 4. Treatment with Lime and Carbon Dioxide; PA1 5. Filtering; and PA1 6. Recovering the inulin by evaporation or precipitation.\nThe inulin may be subjected to heat and/or pH adjustment at some stage in the process to denature inulinase.\nThe biology, chemistry and analysis of inulin and related substances is reviewed in \"Science and Technology of Fructans\", M.Suzuki and N. J. Chatterton, Eds., CRC Press, Boca Raton, Fla. 1993. A review of technology relating to inulin is found in \"Inulin and Inulin-containing Crops\", S. Fuchs, Ed., Elsevier Science Publishers B. V., Amsterdam, 1993. In particular, see Vogel, \"A PROCESS FOR THE PRODUCTION OF INULIN AND ITS HYDROLYSIS PRODUCTS FROM PLANT MATERIAL\", pp. 65-75.\nIt is also known that the maintenance of a balance between the various bacteria that constitute the enterobacterial flora is closely related to man's health and, when pathogenic bacteria which are usually in the minority in the intestines become predominant, the systems of a disease develop. A typical example is diarrhea which occurs as a result of the disappearance of Bifidobacterium bifidum or the action of Escherichia coli or Staphylococcus. Bifidobacterium bifidum is also closely related to infant health. More specifically, Bifidobacterium bifidum is predominant in the enterobacterial flora of healthy infants. It is known that the feces of a dysenteric infant have a markedly lowered Bifidobacterium bifidum content and the relative proportions of the various bacteria in the enterobacterial flora is considerably out of balance.\nOn the basis of these various known facts, a typical well-balanced enterobacterial flora is now considered to be such that Bifidobacterium bifidum, which has excellent staying and proliferation potencies in the intestines, is most predominant at all times.\nIn view of the useful activities of Bifidobacterium bifidum, various kinds of Bifidobacterium bifidum-containing preparations and dairy products have recently been developed for the purpose of increasing the amount of Bifidobacterium bifidum in the intestines.\nIt is considered that the most essential factor for the proliferation of Bifidobacterium in the intestines is saccharides.\nRecent, studies found that fructooligosaccharides are effective in promoting the proliferation of Bifidobacterium bifidum. Fructooligosaccharides (FOS) are natural substances composed primarily of fructose molecules. They belong to a group of carbohydrates that occur in many different plants. FOS are indigestible oligosaccharides that are members of the inulin subclass of fructosans, polymers composed of fructose residues. Specifically, inulins are glucofructosans, carbohydrate polymers consisting of a chain of fructose resides linked by glycosidic bonds. FOS are not hydrolyzed in the small intestine and pass through without being digested, reaching the large intestine where intact they are selectively fermented and utilized many intestinal microorganisms. FOS can be utilized efficiently by lactobacilli and bifidobacteria, species of bacteria that are beneficial for human health (Hidaka et al. \"Fructooligosaccharides: Enzymatic Preparation and Biofunctions\", Journal of Carbohydrate Chemistry 10(4): 509-522, 1991). Selective fermentation of FOS by Bifidobacterium leads to an increase in the presence of these bacteria and to the production of acetic acid and lactic acid as fermentation endproducts, resulting in a lower pH in the digestive tract and providing a means to prevent the overgrowth harmful bacteria like Escherichia coli, Clostridium perfringens and Clostridium difficile. (Hidaka et al., supra.) Hikada et al., also state that fermentation of FOS can also lead to an increase in the presence of short chain fatty acids and the suppression of undesirable microorganisms such as Clostridium perfringens, C. difficile, or E. coli and the toxins they produce. FOS can be utilized most efficiently by bifidobacteria, which are believed to be highly beneficial organisms (Hidaka, et al.), but cannot be utilized by certain undesirable as E. coli and putrefactive bacteria such as Clostridium perfringens or Clostridium difficile. It has also been shown, H. Hidaka et al., Effects of Fructooligosaccharides on Intestinal Flora and Human Health, 5 Bifidobacteria Microflora 37-50 (1986), that administration of fructo-oligosaccharides (FOS) enhances growth of the bifidobacteria population in the intestine, suppresses production of putrefactive factors, improves blood lipid levels in hyperlipidemia patients and provides relief from constipation.\nFructooligosaccharides (FOS) can be produced enzymatically through chemical techniques or by extraction from natural substances. FOS occur in nature in many kinds of plants, including onions, dahlias, garlic, shallots, artichokes, wheat, rye, bananas, asparagus and tomatoes, that are commonly part of a human diet (Speights et al., \"Fructooligosaccharides-A Low Caloric Bulking Agent And More From Sucrose\", Carbohydrates in Industrial Synthesis, ed. M. A. Clarke, Proceedings of the Symposium of the Division of Carbohydrate Chemistry of the American Chemical Society, 1992).\nIt has been shown, A. Hata, The Influence of Neosugar on the Lipid Metabolism of Experimental Animals, Proc. 1.sup.st Neosugar Res. Conference, Tokyo (1982), that fructo-oligosaccharides (FOS) in the diet of experimental animals cause reduction of blood sugar, serum cholesterol, triglycerides, phospholipids; significant improvement in the HDL/LDL ratio; an increase in free fatty acids; and significant decreases in total cholesterol in lipidemia cases.\nAnimal toxicology studies have shown no evidence of toxicity, mutagenicity, or carginogenic effects due to FOS (Clevenger et al., \"Toxicological evaluation of neosugar: genotoxicity, carcinogenicity, and chronic toxicity\", Journal of the American College of Toxicology 7:643-622, 1988). Indigestible oligosaccharides such as FOS can be added to a nutritional product to create an environment in the gastrointestinal tract that is not conducive to the growth of microbial pathogens. Such a nutritional product can also be useful in the prevention of diarrhea caused by these pathogens. FOS is used in Japan in many food products and has been added to infant formula (Fructooligosaccharide Information Package, Coors Bio Tech, Inc. May 1990).\nCertain bacteria have also been shown to be beneficial to human gastrointestinal health. The intestinal flora of the human gut contains some 100.times.10.sup.9 viable bacteria, representing 100 or more different species. The major bacteria of the intestine can be roughly divided into three groups: (a) lactic acid bacteria, including lactobacilli, bifidobacteria, and streptococci; (b) anaerobic bacteria; and (c) aerobic bacteria.\nBacteria of the genus Lactobacillus have been used for several hundred years for treating various illnesses. Lactobacilli found in the human intestinal tract include L. acidophilus, L. casei, L. fermentum, L. salivaroes, L. brevis, L. leichmannii, L. plantarum, and L. cellobiosus. In recent years, L. acidophilus has been shown to be exceptionally useful in treating conditions such as antibiotic-induced imbalances in the gastrointestinal microflora, hypercholesterolemia, vaginal infections, E.coli infection, oral contraceptive failure, depressed immunity, cancerous tumors, chronic granulomatous disease, and lactose indigestion. A. G. Shauss, Method of Action, Clinical Application, and Toxicity Data, 3 J. Advancement Med. 163 (1990). In vitro studies have shown L. acidophilus to have an inhibitory effect on the growth of pathogenic bacteria such as Campylobacterpylori, Staphylococcus aureus, Pseudomonas aeruginosa, and Sarcina lutea. K. M. Shahani et al., Natural Antibiotic Activity of Lactobacillus Acidophilus and Bulgaricus, 11 Cultured Dairy Products J. 14 (1976).\nThe beneficial effect of L. acidophilus is further illustrated by preliminary evidence that L. acidophilus inhibits the toxic activities of bacteria in patients with chronic kidney failure. M. L. Simenhoff et al., Biomodulation of Uremic Pathophysiology in Man, abstract presented at Am. Soc. Of Nephrology Meeting, Baltimore, 1992. Such patients often have toxic levels of amines in their blood due to bacterial overgrowth in the small bowel. Consumption of high levels of freeze dried bacteria drastically reduces levels of these toxic amines. These results demonstrate the ability of L. acidophilus to exert a positive effect on the microflora of the intestines.\nBifidobacteria are also known to exert a beneficial influence on human health. These bacteria exert antimicrobial activity in the human intestine by producing lactic acid and acetic acid as a result of carbohydrate metabolism. These acids lower the intestinal pH, thereby inhibiting overgrowth of gastrointestinal pathogens. Therapeutic applications of bifidobacteria are indicated for the management of diarrhea and constipation, and the management of hepatic encephalopathy with hyperammonemia. Additional benefits include the production of B vitamins and breakdown of carginogenic N-nitrosamines.\nBifidobacteria can be significantly reduced in elderly ispeople due to a reduction of secreted gastric juices. The bifidobacteria population in adults is much more stable, however, changes in the diet, administration of antibiotics, exposure to gamma radiation or x-rays, disease, stress and other disturbances can result in overgrowth of potentially pathogenic bacteria, decrease in beneficial bacteria (lactobacilli and bifidobacteria), and in the gut is associated with various forms of diarrhea, susceptibility to systemic infections, constipation, vague and acute abdominal symptoms, fatigue, dyspepsia, and presence of carcinogenic metabolites. Reestablishment of a normal balance of gastrointestinal flora can be accelerated, and such normal balance maintained, with dietary administration of lactobacilli and/or bifidobacteria such as in the present invention. Even when the Bifidobacteria have become temporarily predominant in intestines, its influence may readily change. Thus, it is preferred to promote the growth of not only Bifidobacteria but also other useful bacteria such as lactic acid bacteria in order to stably obtain the effect of improving intestinal floras. In other words, the influence of lactic acid bacteria and so forth may also be increased together with that of Bifidobacteria.\nLactobacilli and bifidobacteria produce organic acids that reduce intestinal pH and thereby inhibit the growth of acid-sensitive undesirable bacteria. Lactobacilli produce lactic acid, hyrdrogen peroxide, and possibly acetic and benzoic acids. Bifidobacteria produce short chain fatty acids (SCFA) such as acetic, propionic, and butyric acids, as well as lactic and formic acids. The most plentiful short chain fatty acid produced by bifidobacteria is acetic acid, which has a wide range of antimicrobial activities against yeasts, molds and other bacteria. Additionally, short chain fatty acids support normal gastrointestinal function by increasing colonic blood flow, stimulating pancreatic enzyme secretion, promoting sodium and water absorption, and potentiating intestinal mucosal growth. Bifidobacteria are also known to deconjugate bile salts to free bile acids, which are more inhibitory to susceptible bacteria than are the conjugated forms. Further, lactobacilli and bifidobacteria are able to produce other antimicrobial substances, such as bacteriocins, that inhibit the growth and proliferation of harmful bacteria in the gut.\nIn view of the foregoing, it will be appreciated that a composition for alleviating irritable bowel symptoms and for improving and maintaining gastrointestinal health comprising aloe, that has antibacterial and antiviral activity, and soluble dietary fibers such as fructo-oligosaccharides and inulin, that provide the typical advantages of dietary fiber and additionally are low in calories, does not affect blood glucose or insulin levels, further including beneficial/friendly bacteria which favor the growth of other beneficial bacteria in the gastrointestinal tract while at the same time inhibiting the growth of potentially pathogenic or harmful microorganisms would be a significant advancement in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to network elements capable of providing new and enhanced broadband, information-related services to consumers. Specifically, the present invention provides the ability to receive and respond to diverse subscriber requests by transmitting information of virtually any kind. The information transmitted can include movies, video games and other entertainment programs, educational information and programs, scientific and other research-related database information, and the like, and is hereinafter referenced generically as \"information\" or \"titles \".\n2. Description of the Related Art\nConsumer demand for enhanced on-site entertainment and information services is on the rise. Examples of such entertainment and information services include so-called \"on-demand\" video, interactive video games, database research, \"home-shopping\" and the like. Numerous service providers are currently eager to tap this demand, and are therefore expressing interest in schemes for providing such services to consumers. The efforts currently being expended by service providers, as well as the current state of progress in this field of technology are described, e.g., in TIME Magazine, Apr. 12, 1993, at 49ff., and Newsweek, May 31, 1993, at 39ff.\nIn order to be economically practical and viable, a system offering the wide range of information-related services just described should preferably be able to (i) permit access to large volumes of dam (i.e., to a large number of titles), (ii) service a large number of consumers, and (iii) transfer data (i.e., titles) to the end user quickly, preferably in \"real time\". Further, the consensus prefers a system arranged as a centralized facility capable of communicating with remote stations (e.g., homes, offices, public gathering locations, etc.) via existing or next-generation communications media.\nParticular on-going efforts related to the present invention include a system being developed by Broadband Technologies. General elements of the Broadband Technologies system, as understood, are illustrated in FIG. 15. As shown, a subscriber's remote location 151 is connected over a standard telecommunications network 152 to a central facility 153. The central facility 153 awaits service requests from its subscribers.\nWhen a given subscriber wishes to view a selection (e.g., a movie), he or she places a telephone call to the central facility using a standard telephone 154. The call is routed via standard telephone switches 155 (e.g., EO and AT switches) over the telephone voice path of a standard telephone line 156 (e.g., 24-gauge twisted pair copper wire). The telephone call is directed to a Voice Response Unit (VRU) 157 located in the central facility. The VRU forwards the subscriber request to an on-site attendant 158, e.g., through a paper printout. The operator 158 reviews the request, then retrieves the requested selection from an in-house library 159. Finally, the operator 158 locates an available playback unit 170 in a playback battery 171 and loads and plays the selection accordingly.\nThe output signal of the playback unit 170 is supplied to an interface unit 172 coupled between the playback unit 170 and the VRU 157 in the central facility 153. The interface 172 compresses the signal received and forwards the compressed signal via the VRU 157 to the telephone voice path of the telephone line 156. From there, the signal is routed to the subscriber's telephone 154, where a further interface 173 retrieves the signal, decompresses it, and passes it to the subscriber's television set 174.\nAs such, the Broadband Technologies system succeeds in providing an adequate transfer of video data over existing switches and transmission lines. First, the signal received at the remote location 151 corresponds in quality approximately to that of a conventional slow-play VHS mode signal. Furthermore, this system, while not truly \"interactive\", does provide limited system responses to remotely transmitted user requests. However, the system suffers under a range of inefficiencies, limitations on size and flexibility, and limitations or interactiveness, discussed in greater detail below. Finally, the system requires direct linkage to the subscriber's telephone 154, which monopolizes the telephone line and requires electrical connection between the telephone and television.\nAn additional pilot system being developed, known as the AT&T Leesburg Project, overcomes at least some of the problems noted above. General elements of the Leesburg system, as understood, are illustrated in FIG. 16. As apparent from the figure, the Leesburg system is very similar to the Broadband Technologies system, except that the request and response signals are not routed via standard telephone lines. Instead, the request and the response transmission are routed between the remote location 161 and the central facility 163 via a fiber-optic cable system 162.\nIn operation, a subscriber request is initiated via a customized telephone or control unit 164 connected to a special cable converter box 165 installed at the remote location 161. The request is transmitted over the fiber system 162 to a Voice Response Unit (VRU) 166 located at the central facility 163.\nThe operations at the central facility are the same as those described above with respect to the Broadband Technologies system, as indicated by the use of like reference numbers. The output signal of the playback unit 170 is supplied to an interface unit 167 directly coupled to the fiber system 162. The special cable converter box 165 at the remote end retrieves the signal and passes it on to the subscriber's television 168.\nSince the Leesburg system benefits from the larger data transmission capacity of fiber-optics over standard copper wire (approx. 150 Mbits/sec versus approx. 1.5 Mbits/sec), the resulting quality of the data signal is enhanced relative to that of the signal delivered by the Broadband Technologies system. Moreover, this system does not monopolize the subscriber's telephone line. However this system fails to overcome the other problems noted, namely inefficiency, limitations on size and flexibility, and limitations on interactiveness.\nThe problems inherent in these systems, mentioned briefly above, will now be described more particularly.\nThe first class of problems involves inefficiency. First, no system yet proposed is fully automated. The human involvement and interaction required renders such systems costly, fallacy-prone and sluggish in response time. Second, pronounced redundancy between subscribers and playback capacity is required. In other words, the system must dedicate one VCR per concurrent user per average title duration. Further, the system must stock multiple copies of each title in order to cover overlapping requests. Clearly, this redundancy adds tremendously to overall cost and complexity of the system. Moreover, the arrangement is very wasteful, given that certain selections (e.g., first-run movies) are extremely popular for a time and then drop dramatically in popularity.\nThe second class of problems involves limitations on size and flexibility of the overall system. More specifically, given the redundancies discussed above, any system designed in this manner is inherently limited in the number of subscribers that can be accommodated and in the number of selections (titles) that can be offered.\nThe third class of problems relates to interactiveness. In each of the systems described, the user merely requests and views the title selected. As designed, the systems are incapable of providing more versatile or elaborate options that would require greater interaction between the user and the system. Examples of possible options available from a more powerful system include obtaining information to assist the user in selecting a title, selecting and reviewing previews of titles, adjusting playback of titles (e.g., pause, fast-forward, etc.), selecting between original language and dubbed versions of titles, etc. Further, truly interactive services are foreclosed from the described systems. Such interactive services of interest to consumers include on-line database searching, networked remote-station video-game competitions, personalized home shopping, and the like."}
{"text": "Streaming applications are applications that deal with a large amount of data arriving continuously. In processing streaming application data, the data can arrive late, arrive out of order, and the processing can undergo failure conditions. It can be appreciated that tools designed for previous generations of big data applications may not be ideally suited to process and store streaming application data.\nEnabling streaming applications to store large amounts of data from a storage perspective can be challenging. There is a need to determine the proper storage primitive that would ideally be suited for building a new generation of streaming applications in conjunction with existing tools like Apache Flink. In using a Lambda architecture, a developer may use a complex combination of middleware tools that include batch style middleware influenced by platforms like Apache Hadoop and continuous processing tools like Apache Storm, Apache Samza, Apache Kafka and others. Batch style processing can be used to deliver accurate but potentially out of data analysis of data. It can be appreciated that “real-time” processing may deliver faster results but could come at a cost of reduced accuracy. Furthermore, there may be a need for two copies of application logic because the programming models of a speed layer are different than those used in a batch layer.\nConventionally, Lambda architectures may be expensive to develop and expensive to deploy and manage in production. In some implementations, as more application like Internet of Things (“IoT”), require continuous processing, it may not be beneficial to use Lambda architectures and conventional style middleware. Therefore, there exists a need for a simpler approach to Lambda to process streaming application data."}
{"text": "1. Field of the Invention\nThis invention relates to the generation of electric power. More particularly, it is concerned with a unique method for operating an electric power plant, involving the conversion of coal to liquid and gaseous fuels which are subsequently utilized for generating electric power.\n2. Discussion of the Prior Art\nSolid carbonaceous materials are well known for their utility in the generation of power, particularly electrical energy. Generally, coal is combusted with air and the exothermic heat of reaction is used to produce high-pressure steam. The steam in turn is expanded through a turbine to generate mechanical or electrical energy. Similarly, natural gas and other gaseous fuels have been combusted and the heat generated used to form high-pressure steam for the generation of electrical power.\nThe electrical industry has developed highly efficient, large generators driven by expanding steam. However, one problem in the generation of electrical power from steam results from the greatly varying demands for electrical energy. Steam generators are not well suited for producing greatly varying amounts of steam, but rather are designed for base-load or constant-load types of operation. To provide for peak-load and reserve-load demands, the open-cycle gas turbines have generally been utilized because of their quick-startup capability and relatively low capital cost. Open-cycle turbines, however, require special fuels which are noncorrosive to the turbine blades. Generally it has been found to be uneconomical to combust coal or residual oils directly in the combination chamber of a gas turbine, because the fuel contains high amounts of ash and sulfur. Due to the incomplete combustion, such high-ash solid fuels generally produce solid particles which are abrasive and corrosive to the metal turbines. When such particles are entrained in the flue gas that is passed through the turbine, they deposit on the blades and erode the blade surfaces. When this corrosion occurs the blade is damaged, reducing the effiency of the unit, and the passages in the turbine become clogged. Further, the fine particles may deposit down-stream in heat-exchange surfaces and impair thermal efficiency. Similar problems are encountered when burning ash-producing liquid petroleum products.\nPrevious methods in which the fuel gas was cleaned prior to being introduced into the gas turbine were either impractical, unduly costly, or both.\nThe aforesaid problems of base-load and peak-load demand, combined with the special fuel requirements for gas turbines, are substantially avoided by the present invention, which integrates the production of a base-load power generation capability combined with the production of a clean, normally liquid fuel which is particularly useful for the generation of supplemental peak-load power.\nBecause coal and other solid carbonaceous materials often contain sulfur compounds, the combustion of coal for power production can also cause serious air pollution problems. Also, because of the very large volumes of gases produced by combustion, it is very expensive to remove the polluting sulfur compounds after combustion. These sulfur-removal problems and air-pollution problems have led to processes for the gasification of coal to produce a clean fuel gas wherein the sulfur is removed from the fuel prior to combustion. One problem, however, with such gasification processes is that only a low heat value gas is produced, and it is generally not economical to transport a low heat value gas over great distances. This has led to proposals for large on-site or \"mine-mouth\" gasification and power generation plants where the low heat value gas is immediately converted to electrical power for transmission. Such on-site gasification and power generation processes solve the problem of low heat gas transportation and sulfur-removal problems, but are not economical for producing greatly varying amounts of fuel as is needed for peak-load generation of power, either because it is too expensive to store gaseous products for subsequent use in gas turbines or because the capital expense of providing for greatly increasing the gas production rate and gas cleanup rate for peak-load demand is uneconomical.\nThe aforesaid problems are substantially avoided by the subject invention, which process provides for satisfying base-load power demand. The process integrates the production of a combustible fuel gas for meeting base-load power generation requirements with a process for producing easily storable fuels which are substantially free of sulfur and other impurities. These readily liquified fuels are suitable for storage and, when needed, for supplemental peak-load power generation.\nU.S. Pat. No. 3,868,817 discloses a process for the generation of mechanical and electrical power from a purified fuel gas produced from solid carbonaceous fuels. The purified fuel gas is used to generate power using gas turbines.\nAnother patent, U.S. Pat. No. 3,986,349, discloses a process for generating electrical power from solid carbonaceous material in open-cycle gas turbines to meet variable power demands. The process involves the conversion of coal to a combustible synthesis gas or \"syngas\" by reaction with steam and oxygen. The synthesis gas is then divided into two portions, one of which is contacted with Fischer-Tropsch and hydrogenation catalysts to produce synthetic hydrocarbons ranging from methane and ethane to C.sub.22 or higher. The normally gaseous portion of the product is separated and recombined with the second portion of the synthesis gas stream and the combined streams are combusted and utilized to operate an electricity-generating gas turbine. The normally liquid hydrocarbon product (C.sub.5 -C.sub.22) is stored and utilized as fuel for a gas turbine to produce supplemental power for peak-load demand.\nOne of the more recent developments with regard to clean burning coal fired power plants has been the utilization of oxygen fired coal gasifiers coupled with combined-cycle power generators. In this type of plant a clean gaseous fuel obtained from the coal gasifier is burned in a gas turbine and the hot off-gases are used to generate steam. The steam is then used in a steam turbine to generate power. Heat recovered from the gasifier is used to generate additional steam, both for purposes of providing steam needed in the gasifier and as additional feed to the steam turbine.\nOne disadvantage to such plants has been that most coal gasifiers presently available are such that their output is neither easily nor economically varied to match the sometimes dramatic variations in power demand. The same problem results with regard to combined-cycle power plants with low efficiency if output must be reduced by more than about 20%. Many combined-cycle power plants are therefore operated in an on/off control mode.\nA potential solution to this problem would be to use the syngas during times of low power demand to produce storable fuels. The oxygen-blown coal gasifiers have the advantage that they are similar to the gasifiers used for preparation of liquid fuels from coal. There are at present two processes available for use with gasifier-generated syngas for production of storable fuels: co-production of methanol and production of liquid hydrocarbons by the Fischer-Tropsch synthesis. Co-production of methanol suffers from the fact that almost all gasifiers presently contemplated for that purpose yield a synthesis gas with a low ratio of hydrogen to carbon monoxide (e.g. 0.5 to 1.0), which severely limits the amount of methanol obtainable without an additional and expensive water-gas shift. Production of hydrocarbon liquids via the Fischer-Tropsch route has the disadvantage that significant amounts of methane and ethane are produced and these are expensive to condense or otherwise separate. Also, as thermal energy is lost in the production of non-storable products from the syngas, the efficiency of the total process is greatly reduced. Our invention provides a system which overcomes these shortcomings in a way which results in an overall power generating scheme which is superior to anything previously suggested."}
{"text": "A chip resistor as an exemplary prior-art chip component, for example, includes an insulative substrate such as of ceramic, a resistive film formed by screen-printing a material paste on a surface of the insulative substrate, and an electrode connected to the resistive film. In order to adjust the resistance of the chip resistor at a target resistance value, a trimming trench is carved in the resistive film by applying a laser beam to the resistive film through laser trimming (see PTL1).\nAnother exemplary prior-art chip resistor is disclosed in PTL2. The disclosed chip resistor is configured such that a pair of electrodes are provided in spaced relation on a lower surface of a chip-shaped resistor made of a metal. The resistance of the chip resistor cannot be adjusted at a desired resistance value."}
{"text": "Organic EL display devices are known as thin profile, high image quality, and low power consumption display devices. The organic EL display device has formed therein a plurality of pixel circuits arranged in a matrix, the pixel circuits including organic EL elements, which are light-emitting electrooptical elements driven by current, driving transistors, and the like.\nThe organic EL elements have been known for having a decrease in light emitting efficiency due to deterioration over time, resulting in a decrease in light-emitting luminance. FIG. 17 is a drawing for describing the effect that deterioration over time of the organic EL elements has on image display. More specifically, FIG. 17(A) shows a situation in which the same pattern is displayed over a long period of time, and FIG. 17(B) shows a situation in which all pixel circuits are applied a signal for the same luminance after the same pattern was displayed over the long period of time. As shown in FIG. 17(A), the cumulative light-emitting time for organic EL elements (hereinafter, “organic EL elements in a first region PA”) in the pixel circuits in the region PA (hereinafter, the “first region”) where bright display is performed over a long period time is longer than that of organic EL elements (hereinafter, “organic EL elements in a second region PB”) in pixel circuits in the region PB (hereinafter, the “second region”) where dark display is performed over a long period of time. Thus, the organic EL elements in the first region PA undergo a decrease in light-emitting efficiency due to greater deterioration than those in the second region PB. As a result, as shown in FIG. 17B), so-called screen burn-in occurs in the first region PA. Specifically, display of the same luminance as the second region PB normally should occur in the first region PA, but display of a lower luminance than the second region PB occurs in the first region PA.\nFIG. 18 is a drawing for describing the decrease in luminance of the organic EL elements. Here, a fixed current is assumed to be fed to the organic EL elements. As deterioration over time of the organic EL elements progresses, impedance increases in the organic EL elements. As a result, as shown in FIG. 18, forward-biased voltage applied to the organic EL elements increases as deterioration over time of the organic EL elements progresses. As described above, light-emitting efficiency decreases as deterioration over time of the organic EL elements progresses, and as a result, the decrease in luminance occurs as shown in FIG. 18. The deterioration over time of the organic EL elements in the second region PB has not progressed as much as those in the first region PA, and thus, there is not as much decrease in luminance in the second region PB. On the other hand, the deterioration over time of the organic EL elements in the first region PA has progressed more than in the second region PB, and thus, there is a greater decrease in luminance in the first region PA. As a result, the display state shown in FIG. 17(B) occurs.\nIn relation to the present invention, Patent Document 1 discloses a pixel circuit that compensates for increase in forward bias voltage resulting from deterioration over time of organic EL elements. FIG. 19 is a circuit diagram showing a configuration of the pixel circuit 91 disclosed in Patent Document 1. In FIG. 19, for ease of explanation, the reference characters in the drawing of Patent Document 1 are modified. The pixel circuit 91 has one organic EL element OLED, six transistors T11 to T16, two capacitors C11 and C12, and a variable bias voltage source VS. The transistor T12 is of a p-channel type, and the transistors T11 and T13 to T16 are of an n-channel type.\nFirst, a scan wiring line Sj is selected and the transistor T11 turns ON, and a voltage based on the data signal fed from a data wiring line Di is written to the capacitor C11. Next, the selection of the scan wiring line Sj ends and the transistor T11 turns OFF, and control lines Vg13j and Vg15j are selected. As a result, the transistor T13 turns ON, and a drive current based on a voltage between source and gate of the transistor T12 is fed to the organic EL element OLED. Also, the transistor T15 turns ON, and the gate potential of the transistor T16 becomes equal to the anode potential of the organic EL element OLED based on the drive current. The anode potential Pi of the organic EL element OLED changes due to deterioration of the organic EL element OLED. Here, using the variable bias voltage source VS, a source potential Ps of the transistor T16 is set according to the following formula (1).Ps=Pi−Vth  (1)\nHere, Vth refers to the threshold voltage of the transistor T16.\nBy setting the source potential Ps of the transistor T16 according to formula (1), it is possible to extract the increase in forward bias voltage resulting from the deterioration of the organic EL element OLED as the source/drain current of the transistor T16. After the source/drain voltage of the transistor T16 is determined, the selection of the control line Vg15j ends and the transistor T15 turns OFF, and then the control line Vg14j is selected and the transistor T14 turns ON. Thus, the potential of the gate terminal of the transistor T16 decreases based on the source/drain current of the transistor T12. As a result, it is possible to perform luminance compensation based on the increase in forward bias voltage resulting from deterioration over time of the organic EL element OLED. Therefore, it is possible to mitigate a decrease in luminance of emitted light resulting from the deterioration over time of the organic EL elements OLED."}
{"text": "1. Field of the Invention\nThis invention relates in general to methods and apparatus for collecting and testing a viscous liquid or gel. Apparatus and methods are disclosed for collecting and pooling the contents of a plurality of relatively small containers. The invention is especially useful for microbiological testing a product sample with respect to sterility, bacteriostasis, or preservative effectiveness.\n2. Description of the Related Art\nIn the production of consumer products in the health care or pharmaceutical industries, in which liquids of various viscosities are sold in relatively small containers for individual use, microbiological testing on a statistical basis may be required to ensure that the products are substantially free of microorganisms such as bacteria, fungi and molds. Recognized procedures for such microbiological testing may involve the direct transfer method or the membrane filter method. These methods are described in the United States Pharmacopeia (USP) 23 (1994), U.S. Food and Drug Administration Guidelines (FDA), and the Japanese Pharmacopeia XII suppl. II, all incorporated herein by reference.\nOne conventional method of microbiological testing is the antimicrobial preservative-effectiveness procedure (or \"preservative efficacy\" procedure), normally a direct transfer method. In the direct transfer method, microorganisms are inoculated directly into a sample solution, provided that the sample size of that solution is of sufficiently ample size (normally at least 20 mL). The solutions are then plated out in conventional media and subsequently examined at periodic intervals of 7, 14, 21 and 28 days. When the amount of sample obtained from a single product container is too small to allow the sample to be directly inoculated, a plurality of containers from a sample lot may be pooled together (under sterile conditions to avoid contamination) and then the pooled contents directly inoculated.\nAnother standard microbiological test, sterility testing, typically involves the membrane filter method. In the membrane filter method, a test sample is filtered through a membrane filter. Typically, the filter is then rinsed to remove any microbiocides, after which the filter holding the retentate from the test sample, or a preselected fraction of this filter, is incubated in a suitable medium. Optionally, the medium may be added onto the membrane filter while remaining in the filtration device.\nIn the health care industry, in order to obtain an adequate amount of test solution for sterility testing a product, twenty bottles or more per sample lot may need to be pooled or combined. In order to obtain an adequate amount of test solution for the previously mentioned preservative efficacy test, as many as eighty containers or bottles per sample lot may need to be combined. The pooling of such large numbers of bottles can be very time consuming and difficult when the sample containers are small, the contents are relatively viscous, and/or the bottle is designed to expel the liquid contents dropwise. If the containers have a separate top piece or dispenser tip that is removable from the container, then the tip can be removed, the bottle inverted, and the solution dispensed relatively easily, especially when the solution is non-viscous. Typically, however, the dispensing tip snuggly fits into the neck portion of a bottle, and removal of such container tips can be time consuming and laborious.\nIt may be particuarly problematic or difficult to transfer a viscous solution from a single-piece container having a dispensing tip that is an irremovable or integral part of the container holding the solution. Such containers include form-filled bottles or tubes, either unit-dose or multi-dose containers. A conventional method of collecting a sample from a container having an irremovable tip involves repetitively squeezing the outside surface of the container to force its contents out its dispensing orifice as quickly as possible. To dispense the entire contents of the necessary number of bottles by such a method may be a strenuous and monotonous task, particularly when the bottles are semi-rigid and designed to deliver droplets. Also, such a method may result in sample contamination as a result of air being drawn back into a container.\nThere are many devices available for sterility testing of liquids. For instance, U.S. Pat. No. 4,036,698 relates in general to the testing of pharmacological products. The patent describes and claims an apparatus designed to carry out a membrane filtration method wherein the solution to be tested is added to a filter apparatus comprising a cylindrical canister having two ports at one end and a single port at the opposite end. One of the two ports at the one end is provided with a hydrophobic microporous filter that allows the flow of air therethrough in either direction while screening out any microorganisms. A second membrane filter is positioned within the cylinder such that test solution passes therethrough, trapping bacteria or fungi on the membrane filter. This apparatus does not solve or address the above-mentioned difficulties of collecting a sufficient amount of sample contents from many relatively small containers, particularly containers designed for the delivery of droplets.\nIn view of the above, there is a need for a more efficient method and apparatus for collecting samples for the microbiological testing of products produced in the pharmaceutical or health care industries. There is a need for a more efficient method of collecting samples from product containers designed to deliver droplets during individual use."}
{"text": "Telecommunication networks provide for the transmission of information across some distance through terrestrial, wireless or satellite communication networks. Such communications may involve voice, data or multimedia information, among others. One particular example of transmission of data or multimedia information involves a content delivery network (CDN). CDNs are increasingly used to distribute content, such as videos, multimedia, images, audio files, documents, software, and other electronic resources, to end users on behalf of one or more content providers. Using a CDN allows the content providers to increase the speed and reliability of content delivery without deploying additional infrastructure. Moreover, a CDN allows for the distribution of the content through one or more existing networks without the need to store the content within the existing networks.\nTypically, a CDN includes several content servers from which the content can be supplied to a requesting end user. In one example, these content servers may be accessed through a telecommunications network to which the end user is in communication. The network may include any number of components to facilitate the connection of the end user to the requested content, such as routers, Internet Service Provider networks, other intermediate networks, and the like. In general, the content available from the CDN is stored on one or more edge clusters connected to the CDN or other upstream content providers. Requests for content are then transmitted by the CDN to the edge clusters or content providers to provide the content to the requesting customers. However, the CDN may desire to direct the end user's computing device to a specific content storage device or server.\nIt is with these observations in mind, among others, that various aspects of the present disclosure were conceived and developed."}
{"text": "1. Field of the Invention\nThe present invention generally relates to storage containers for storing foods and other products that bear liquids which are likely to be exuded from the products. More particularly, the present invention relates to food storage containers incorporating a reservoir which contains absorbent material for collecting and absorbing excess liquids which exude from food, etc., in the form of biofluids, marinades, saline, juices, water, etc., so that the liquids and their associated microorganisms are held separate from the food while contained in the storage containers.\n2. Description of the Related Art\nExcess moisture within food storage containers can cause premature spoilage of food products which are stored in the containers because the moisture provides a favorable environment for the growth and reproduction of microorganisms. Excess moisture in a storage container also can lead to leakage of fluids from the storage container which can cause contamination of other foods and items about the container. Attempts at controlling excess moisture in food storage containers, such as trays (rigid and flexible) and bags, have met with some success. These prior art devices include: (1) pre-formed trays configured for the insertion of absorbent pads or absorbent sheets with the food products resting on the absorbent pad; (2) trays with built-in reservoirs arranged to trap excess moisture exuded from the products, with some of the trays including a cover which allows fluids to drain from the product into the reservoir but which partially restricts the fluids from reemerging past the cover following shaking or movement of the tray; and (3) trays or packs made from multiple layers of material with one layer being liquid impervious, the second layer being formed with perforations to allow fluids to enter, and an absorbent media sandwiched between the two layers to absorb and retain the entering fluids.\nSeveral aspects of the prior art food containers with absorbent media limit the overall effectiveness of the containers and thereby fail to adequately address the need for absorbing excess liquids exuded from food products stored in containers. Particularly, a pad of absorbent material, typically formed of paper or a fluff-pulp structure, absorbs fluids but can not retain a great quantity of fluids when the pad is under pressure, such as when food is placed directly upon the pad or when food is placed on an upper layer of a tray which incorporates a pad sandwiched below the upper layer. In these configurations, the weight of the food upon the upper layer or directly on the absorbent pad compresses the pad and reduces the capacity of the pad to absorb liquids. Additionally, the amount of paper or fluff that would be required to produce a pad with sufficient absorbency to absorb and retain high levels of liquid is likely to require the pad to be too large for practical use. Furthermore, the prior art food storage trays usually have not incorporated reservoirs of a sufficient size to hold high levels of excess fluid without physically weakening the trays, thereby increasing the tendency of the trays to crack and leak fluids.\nTherefore, there exists a need for devices and methods that address these and other shortcomings of the prior art.\nBriefly stated, the present invention relates to storage containers for storing foods and other products that bear liquids which are likely to be exuded from the products. In a preferred embodiment, the storage container includes a base that defines an interior, and a lid that defines a reservoir. The lid is configured to engage the base so that food products placed within the interior may be retained therein. An absorbent material is arranged within the reservoir of the lid and is retained therein by a cover formed of liquid permeable material. The cover permits liquid exuded from the food products to enter the reservoir so that the absorbent material may absorb the liquid, such as when the storage container is inverted.\nIn some embodiments, the storage container includes a base defining an interior which is adapted to receive food products. Additionally, means for engaging the base and means for absorbing liquid exuded from the food products are provided.\nOther embodiments of the present invention may be construed as providing methods for storing food products. In a preferred embodiment, the method includes the steps of: providing a storage container having a base and a lid, the base defining an interior; placing food products within the interior of the storage container when the base is oriented in a substantially upright configuration; engaging the base with the lid such that the food products are retained within the storage container; inverting the storage container such that the lid is disposed substantially beneath the base; and enabling liquid from the interior of the storage container to permeate into the lid such that the liquid is substantially removed from the interior and retained by the lid.\nOther features and advantages of the present invention will become apparent upon reading the following specification, when taken in conjunction with the accompanying drawings."}
{"text": "Communications networks typically include a plurality of data terminal equipment (DTE), also referred to as terminals or stations. A computer connected to a network is considered to be a terminal. Terminals communicate with each other through a transmission medium which may, for example, be a coaxial cable in the case of Ethernet networks (which will be understood to include Cheapernet networks in the discussion below).\nThe various terminals of a network send information messages and receive those sent by other terminals. An information message is composed of a set of elementary information blocks, called frames, that contain a certain quantity of binary information. Each frame is structured and contains data which define its beginning and its end, the address of the terminal for which the message is intended, the address of the send terminal, the data length, and other useful data.\nA local area network (LAN) is a network limited to a geographic location within a restricted area, the distances between the various stations being on the order of a few meters to tens of meters, even reaching a few kilometers. LANs of the CSMA/CD type are frequently used in current practice. They are standardized by the Institute of Electrical and Electronic Engineers (IEEE) in accordance with a standard known as 802.3 (adopted by the International Standards Organization (ISO), as ISO Standard 8 802.3).\nThis standard defines a number of provisions that establish the communication mode between the various terminals as well as the frame format and the protocol governing dialogs between the various stations. A protocol defines the access rules to the various stations.\nEthernet networks have a data transmission speed of 10 Mbits/s and their standardized transmission medium is a coaxial cable with a characteristic impedance of 50 ohms.\nIt is known that a computer is composed of one or more central processing units, input/output processors, random access memories, and read only memories associated with all these processors, input/output controllers, as well as several peripherals such as disk memories or input/output peripherals allowing exchange of data with users (screen terminals, printers, etc.), these peripherals being associated with respective peripheral controllers.\nAll the component parts listed above (aside from the peripherals) are disposed on a set of boards whose dimensions are standardized. These boards are generally connected to a parallel bus which provides communications between the various processors and data transport between the various boards, as well as providing electrical power.\nOne bus commonly used in present-day practice is called MULTIBUS II, a registered trademark of the INTEL COMPANY. The architecture of such a bus is structured around a principal bus of the parallel type standardized according to IEEE Standard 1296.\nThe rapid technological evolution of networks and the increasing numbers of terminals in use are leading to the development within the computers themselves of programmed communications servers whose role is to reduce the load on the computer's central processing unit by performing some of the management of the messages sent by the various component elements of the computer to the telecommunications network to which it is connected, as well as the messages coming from other terminals on the network. In current practice, such a communications server is built around a microprocessor connected to one or more memories, which microprocessor works with a basic program (simpler than that of the central processing unit) containing specialized modules allowing the bus common to the various component elements of the computer and the network transmission line to be managed. An example of such a program, called communications software, is the program called CNS which is used in the Bull S. A. Company's DN-7XXX series, and also in the CNS-A.sub.0 and CNS-A.sub.1 products in the DPS-7000 computer series.\nExtremely rapid development in the power and processing capacity of central processing units makes it possible for an increasingly large number of input/output peripherals to communicate with central processing units. These input/output peripherals are synchronous or asynchronous peripherals that communicate with their environment via transmission lines having speeds that range from 300 bits/s to 64 kbits/s. These input/output computer terminals communicate with the outside via modems whose function is to adapt the electrical signal delivered by the input/output terminal to the transmission medium connecting the terminal to its outside environment. These modems are, for example, defined by the V-24, V-28, V-11, V-35, and V-36 standards of the Consultative Committee in International Telegraphy and Telephony (CCITT). These various standards also define the transmission modes and protocols for the corresponding transmission links. These are supported physically by transmission media which are in fact cables or sets of individual wires.\nThe large computer systems of today, composed of a computer, a communications server, and the various input/output terminals connected thereto, can have as many as several hundred terminals which are thus connected by the same number of cables or sets of wires to the computer. In practice, it is convenient to connect each computer system to other systems and to the input/output terminals via an Ethernet network. When the terminals are distributed in a geographic location within a restricted area, the existence of a large number of cables and wires poses mechanical and physical problems, such as crosstalk, also occupying too much space."}
{"text": "Anal fissure (or fissure-in-ano), anal ulcer, acute hemorrhoidal disease, and levator spasm (proctalgia fugax) are common, benign conditions of the anal canal which affect men and women. An anal fissure or ulcer is a tear or ulcer of the mucosa or lining tissue of the distal anal canal. An anal fissure/ulcer can be associated with other systemic or local diseases, but it is more frequently present as an isolated finding. The typical, idiopathic fissure or ulcer is confined to the anal mucosa, and usually lies in the posterior midline, distal to the dentate line. The person with an anal fissure or ulcer suffers from anal pain and bleeding, more pronounced during and after bowel movements.\nHemorrhoids are specialized vascular areas lying subjacent to the anal mucosa. Symptomatic hemorrhoidal disease is manifest by bleeding, thrombosis or prolapse of the hemorrhoidal tissues. Men and women are affected. Most commonly, internal hemorrhoidal tissue bulges into the anal canal during defecation causing bleeding. As the tissue enlarges, prolapse pain, thrombosis, and bleeding can ensue. Thrombosis of internal or external hemorrhoids is another cause of pain and bleeding.\nLevator spasm (or proctalgia fugax) is a condition of unknown etiology affecting women more frequently than men. This syndrome is characterized by spasticity of the levator ani muscle, a portion of the anal sphincter complex. The patient suffering from levator spasm complains of severe, episodic rectal pain. Physical exam may reveal spasm of the puborectalis muscle. Pain may be reproduced by direct pressure on this muscle. Bleeding is not associated with this condition.\nThe underlying causes of these problems are poorly understood. However, all of these disorders are associated with a relative or absolute degree of anal sphincter hypertonicity. In the case of anal fissure/ulcer the abnormality appears to be an as yet unidentified problem of the internal and sphincter muscle. The internal sphincter is a specialized, involuntary muscle arising from the inner circular muscular layer of the rectum. Intra-anal pressure measurements obtained from people suffering from typical anal fissure/ulcer disease show an exaggerated pressure response to a variety of stimuli. The abnormally high intra-anal pressure is generated by the internal sphincter muscle. The abnormally elevated intra-anal pressure is responsible for non-healing of the fissure/ulcer and the associated pain.\nAn abnormal pressure response in the anal canal has also been observed in people suffering from symptomatic hemorrhoidal disease. Elevated intra-anal pressures may be a major etiologic factor in the development of this condition. It is postulated that the pain associated with acute hemorrhoidal disease is caused in part by spasm of the internal anal sphincter muscle. Similarly, the pain associated with levator spasm is induced by the muscle spasm itself.\nVarious therapies have been devised to treat these problems. Typical, non-surgical therapy includes bulk laxatives and sitz baths. Sitz baths are helpful because they induce relaxation of the anal sphincter mechanism. Topical anal therapy is used to promote healing, relieve pain, and reduce swelling and inflammation. Many preparations have been tried, including those containing local anesthetics, corticosteroids, astringents, and other agents. None of these preparations adequately addresses the underlying problem of sphincter spasm. Consequently, none has been show conclusively to favorably alter the time course to healing or to reliably ameliorate associated pain.\nThose cases of anal fissure/ulcer or hemorrhoids recalcitrant to medical therapy are often referred for surgical treatment. In keeping with the proposed etiology of anal fissure/ulcer, the current standard surgical procedure for treatment of anal fissure is lateral internal anal sphincterotomy. In this procedure, the internal anal sphincter muscle is partially cut, thereby reducing the intra-anal pressure. The lowered pressure allows the fissure/ulcer to heal and also relieves the associated pain. Surgical hemorrhoidectomy removes the redundant hemorrhoidal tissue. Many surgeons will perform concomitant limited internal anal sphincterotomy to lower anal canal pressure. There is no successful surgical treatment for levator spasm.\nOver the past five years a third component of the autonomic nervous system, the enteric nervous system (ENS), has been described and elucidated. This neural network innervates the gut continuously from esophagus to anus. It is composed of enteric neurons and the processes of extrinsic efferent and afferent neurons of the traditional autonomic system. This system regulates the motor and secretory function of the gut.\nThe most remarkably feature of the ENS is the diversity of chemical messengers that enteric neurons contain and release. In addition to acetylcholine and norepinephrine, various peptide and non-peptide substances have been identified which appear to function as neurotransmitters. Most recently, nitric oxide (NO) has been identified as an inhibitory transmitter to muscle. Rattan, Chakder, O'Kelly, and others have shown that NO mediates the anorectal inhibitory reflex in animals and man. See, Rattan et al., S. Nitric oxide pathway in rectonal inhibitory reflex of opossum internal anal sphincter, Gastroenterology 103:43-50, 1992; Chakder et al., Release of nitric oxide by activation of nonadrenergic noncholinergic neurons of internal anal sphincter, Am. J. Physiol. 264:G7-12, 1993; and O'Kelly, et al., Nerve mediated relaxation of the human internal anal sphincter: The role of nitric oxide. Gut 34:689-693, 1994, each of which is incorporated herein by reference.\nOrganic nitrates such as nitroglycerin (the trinitrate, NTG), isosorbide dinitrate, isosorbide mononitrate, erythrityl tetranitrate, and others have been used for decades in the clinical setting of angina pectoris. These agents act as physiologic nitric oxide donors. The use of organic nitrates has not been previously proposed for the treatment of anal disease.\nCorticosteroids such as hydrocortisone, have been used for various benign anal disorders for many years. Studies of the effectiveness of this treatment have shown some benefit, but not in a reproducible or significant fashion. It has not been heretofore known to use hydrocortisone in combination with organic nitrates for treatment of anal diseases.\nTopical anesthetics such as dibucaine, lidocaine, pramoxine, and others have been used for treatment of anal pain. It has not been heretofore known to use topical anesthetics in combination with organic nitrates for treatment of anal diseases."}
{"text": "Numerous products today are shipped and sold in reclosable plastic bags. Examples of plastic bags having a reclose feature can be found in U.S. Pat. No. 4,617,683 issued to Christoff on Oct. 14, 1986, and U.S. Pat. No. 4,986,673 issued to Bell on Jan. 22, 1991. When using one of the bags disclosed in the above U.S. patents, the consumer first removes the sealed portion of the bag located directly above the reclosing feature. The consumer then pulls apart or separates the male and female securement members of the reclose feature. The consumer can now remove the desired portion of the product from the bag. The consumer then manually engages the male and female securement members of the reclose feature, thereby temporarily securing the bag and its contents until the next use. However, as is often the case with a bag having a reclose feature, there is no biasing means to keep the bag mouth open after the reclose feature has been pulled apart by the consumer. The lack of a biasing means to keep the package in a fully open condition can be cumbersome and annoying to a consumer who desires to have the bag mouth remain open after the reclose feature has been pulled apart, thereby facilitating access to the product within the bag. Improved access to the product is particularly desirable when the bag contains certain products, many of which are not entirely consumed upon initial opening of the pouch, e.g., coffee, salted snacks, fabric softeners, cleaning products, detergents, pharmaceuticals, etc. Improved access is also highly desirable when the pouch is to be used as a vessel for the reconstituting fluid, e.g., concentrated beverages which are to be diluted with water.\nCommonly assigned U.S. Pat. No. 4,898,477 issued to Cox et al. on Feb. 6, 1990 and which is hereby incorporated herein by reference discloses a self-expanding flexible pouch that can be used as a measuring device for reconstituting a concentrated product contained therein. The flexible pouch includes an extensible stay located in the throat area of the pouch. When the pouch's top portion is removed, the extensible stay expands to a circular configuration opening the pouch's throat area. The consumer can then \"milk\" the concentrated product out of the pouch. The opening of the pouch's throat area allows the consumer then to reconstitute the fluid by adding water to the pouch and then pouring and mixing with the concentrated substance previously \"milked\" from the pouch. The Cox et al. pouch has the advantage of self-expansion, thereby improving access to the product within the pouch. However, this pouch would not be particularly well suited to package a product that required the pouch to be sealed between repeated multiple openings.\nWhen using certain products that require repeated or intermittent use prior to exhaustion of the pouch's contents or when using products which are to be diluted using the empty pouch as a vessel for the reconstituting fluid, the consumer not only desires that the bag be reclosable between uses, but also that the bag mouth remain open. By way of example, liquid fabric softener, granular detergents, and salted snacks are products that the consumer would like to reclose between uses to prevent spilling or preserve freshness, and would also prefer to have a self-opening feature to improve or facilitate product access.\nAccordingly, it is an object of the present invention to provide a flexible pouch having a self-opening and reclosing device suitable for use with a wide assortment of products;\nIt is further an object of the present invention to provide a flexible pouch that will assume a substantially predetermined configuration upon opening;\nIt is further an object of the present invention to provide a flexible pouch that will collapse when discarded and thereby reduce the volume of generated waste material;\nIt is further an object of the present invention to provide a flexible pouch that is sufficiently rigid and stable when opened to allow a consumer to firmly grasp and pour a fluid from the pouch with confidence that the pouch will not collapse;\nIt is further an object of the present invention to provide, in a particularly preferred embodiment, a flexible pouch that can be readily opened without having to use an implement such as a pair of scissors or a knife;\nIt is still another object of the present invention to provide a flexible pouch that accomplishes the aforementioned objectives at minimal cost."}
{"text": "1. Field of the Invention\nThe invention relates generally to the field of rotating machines, such as motors and generators. More particularly, it relates to bearing assemblies used with such machines and devices useful when removing the bearing assemblies.\n2. Description of Related Art\nRolling element bearings, also known as anti-friction bearings, are widely used on many types of rotating machinery and have a limited life span. That life span can be reduced by misuse and adverse operating conditions. Bearing removal is frequently required in the field as part of the rotating machinery's maintenance.\nTraditionally, bearing assemblies have been designed such that the removal of the bearing requires the removal of the entire bearing bracket that houses the bearing and that is mounted to the frame of the rotating machine (e.g., a motor or a generator). FIGS. 1 and 2 show an example of such a prior art bearing assembly.\nBearing 10 was positioned, or housed, in bearing bracket 20. Inner race 12 of bearing 10 contacted shaft 30 of the rotating machine (the body of which is not shown). Outer race 14 of bearing 10 contacted machined cavity 22 of bearing bracket 20. Rolling elements 16 (balls in this case) were positioned between, and in direct contact with, the two races. An inner cap 28 was bolted to bearing bracket 20 with bolts 25 positioned through bolt holes 27. The bearing bracket had a machined rim 24 near its outer edge that fit with very little radial clearance into the machined frame of the rotating machine. That tight fit between the frame and bearing bracket 20 is commonly known as a spigot fit, and bearing bracket 20's rim 24, which fit tightly inside of the machined frame, is known as a spigot. Mounting bolts (not shown, but see FIG. 1, which shows mounting holes 32) secured bracket 20 to the frame of the rotating machine to ensure that spigot 26 remained in contact with the frame. During normal operation, the rotor weight was transmitted through the bearings to the bearing brackets (there is typically one at each end of the machine) and through each bracket's spigot to the frame of the rotating machine, which was supported by the site foundation.\nBearing 10 has been removed from rotating machines, such as motors and generators, using jacking bolts. Jacking holes 34 were positioned in bearing bracket 20. Jacking bolts are normal bolts that were screwed into jacking holes 34, which were located evenly around the outside of the bearing bracket such that they aligned with flat, solid steel surfaces on the frame of the rotating machine. As the jacking bolts were placed in jacking holes 34 and tightened, the bottom of the bolts came into contact with the steel surface of the frame. Gradual and even tightening of all the jacking bolts caused the jacking bolts to pull bearing bracket 20 out of the frame. Bearing 10 typically stayed in place on shaft 30 because inner cap 28 was first un-bolted from bearing bracket 20. Because 10 normally stayed in place on shaft 30, bearing bracket 20 could only be removed a short distance before it had to be supported to keep it from falling onto and possibly damaging shaft 30. The weight of bearing bracket 20 (potentially over 200 pounds for a 300 horsepower (HP) rotating machine) generally required a crane or other large lifting device to be used for bearing removal on 300 HP (and larger) motors.\nTo summarize, the following steps have been taken during removal of the bearing shown in FIGS. 1 and 2. The end of shaft 30 has been supported using any suitable mechanism, such as a crane or jacks. The mounting bolts positioned in mounting holes 32 that held bearing bracket 20 to the frame of the rotating machine have been removed, as have bolts 25 positioned in bolt holes 27 that held inner cap 28 to bearing bracket 20. Jacking bolts have been inserted into jacking holes 34. The jacking bolts have been tightened gradually and evenly until spigot 26 of bearing bracket 20 has been completely pulled out of the frame. The supported shaft end has been lowered down until the rotor core (not shown) surrounding shaft 30 contacts the stator core (not shown) surrounding the rotor core. The shaft support has been removed. A crane has been attached to bearing bracket 20 and removed bearing bracket 20. A traditional bearing puller has been used to remove bearing 20 from around shaft 30.\nAn advance was made with the creation of the bearing assembly shown in FIGS. 3 and 4. Bearing assembly 40 (sometimes called a bearing cartridge) is positioned in a bearing assembly opening 42 of bearing bracket 20. Bearing 10 is seated inside of bearing assembly 40 instead of directly in bearing bracket 20, as shown in FIG. 2. As a result, the end user only had to remove bearing assembly 40—not bearing bracket 20—in order to remove bearing 10. Bearing assembly 40 is considerably lighter than bearing bracket 20 and eliminates the need for a crane.\nBearing assembly 40 includes a cylindrical insert 44, an inner cap 28, an outer cap 46, and bearing 10. Insert 44 was machined on its inside to hold bearing 10 tightly and had a spigot fit with bearing bracket 20. Insert 44 also had mounting bolt holes 48 through which bolts were placed to secure insert 44 to bearing bracket 20. Insert 44 also had insert jacking holes 49 that allowed for removal of the bearing assembly from bearing bracket 20. Insert 44, outer cap 46, and inner cap 28 each had aligned openings (designated generally at 50) that allowed them to be bolted together with bolts 52. Outer cap 46 was a disk that bolted onto insert 44 and that had a small clearance fit with shaft 30 to hold the bearing grease inside of bearing assembly 40. Inner cap 28 also was a disk that had a small clearance fit with shaft 30 for grease retention, but also had two ring-shaped projections 51 and 53, positioned to align with inner race 12 and outer race 14, respectively, of bearing 10. Those projections allowed inner cap 28 to act as a bearing puller.\nDuring bearing removal, the jacking bolts positioned in insert jacking holes 49 pulled bearing assembly 40 out of bearing bracket 20. Because inner cap 28 was bolted insert 44, ring-shaped projections 51 and 53 contacted bearing 10's inner race 12 and outer race 14, respectively, and pulled bearing 10 off of shaft 30 during removal. This eliminated the need for the end user to use a bearing puller. Also, the lighter weight of bearing assembly 40 (versus bearing bracket 20) and the presence of bearing 10 in bearing assembly 40 reduced the risk of the bearing assembly falling sizable distances and damaging the shaft as the bearing assembly was pulled off of the end of the shaft.\nTo summarize, the following steps have been taken to remove bearing assembly 40. The end of shaft 30 has been supported. The mounting bolts running through mounting bolt holes 48 have been removed. Those same bolts have then been used as jacking bolts, inserted into insert jacking holes 49 and tightened gradually and evenly until bearing assembly 40 was removed from bearing assembly opening 42 of bearing bracket 40. Shaft 30 was lowered until the rotor core (not shown) surrounding the shaft contacted the stator core (not shown) surrounding the rotor core. The shaft support was removed. The loose bearing assembly 40 was removed by sliding it off of the end of the shaft.\nThe inventors have discovered certain shortcomings associated with the removal of bearing assembly 40 shown in FIG. 4. While removing bearing assembly 40 from bearing assembly opening 42, a large friction force is induced between bearing 10 and shaft 30 tends to move shaft 30 in the same direction as bearing 10. This shaft motion or force is transferred to, and undertaken by, the bearing or coupling at the other end of the shaft. If that bearing (i.e., at the other end) is a rolling element bearing, the friction force between that bearing and the shaft will cause the entire rotor assembly to be moved in the direction of bearing 10. Alternatively, the load will be undertaken by the rolling elements of the bearing or the inner cap located at the other end of shaft 30."}
{"text": "Various catalysts are known to promote the \"ene reaction\" of olefins with alpha-, beta-unsaturated carboxylic acid esters for the production of unsaturated carboxylic acid esters. For instance, U.S. Pat. No. 3,783,136 and German Offenlengungsschrift both describe the use of AlCl.sub.3 and AlB.sub.3 as catalysts for such reactions. U.S. Pat. No. 4,506,095 describes the reaction of linear alpha-olefins with alkyl acrylates catalyzed by an organometallic catalyst of the formula R.sub.n --Al--X.sub.3-n, wherein R is an organic radical containing between about 1 and 12 carbon atoms, n is the integer 1 or 2, and X is chlorine or bromine. U.S. Pat. No. 3,641,120 describes the reaction of an ester with an olefin in the presence of a combination of a manganic carboxylic acid salt or oxide with a zirconyl carboxylic acid salt or zirconium oxide. The publication by B. R. Snider on J. Org. Chem., vol. 39, no. 2 (1972), p. 255, refers generally to Lewis acid catalysts for ene reactions, but illustrates only the use of aluminum chloride and zinc bromide. U.S. Pat. No. 3,855,255 describes the reaction of carboxylic acid esters by reacting diolefins with methacrylate esters in the presence of an organometallic complex of zero-valent nickel and an electron donor. U.S. Pat. No. 2,093,695 discloses preparation of carboxylic acid esters by reaction of acyloxy compounds with olefinic hydrocarbons catalyzed by activated charcoal, inorganic acids, the halogens and various halides of calcium, boron, cadmium, zinc, calcium and potassium. Co-pending Application Ser. No. 510,309, filed Apr. 17, 1990, U.S. Pat. No. 5,066,829 described an ene reaction catalyzed by tantalum pentachoride. Akermark et al (J. Org. Chem., vol. 43, no. 22 (1978), p. 4387) have reported that the eutectic mixture of AlCl.sub.3, NaCl, and KCl is a superior ene reaction catalyst. U.S. Pat. No. 3,892,788 teaches a ligand-stabilized Pt(II) dihalide complex combined with a Group IVB metal halide as a catalyst for such ene reactions. South African Patent 496,265 describes one such reaction catalyzed by various organo-metal compounds. According to German Offenlegungschrift 2044159, a reaction between acrylic acid esters and dienes is catalyzed by an organometalic complex of zero-valent iron and a triaryl compound of an element of Group V. U.S. Pat. Nos. 3,755,386 and 4,144,257 and Belgian published application 739,208 describe similar reactions using complexes of Group VIII compounds, as well as various compounds and complexes of iron, nickel and cobalt. U.S. Pat. Nos. 4,009,203 and 3,534,087, German Offenlegungsschrift 3149979 and World Pat. No. 8100846 describe related reactions of acids and olefins catalyzed by and acyloxy-stannic trihalide or a perfluorosulfonic acid resin or a crystalline metal silicate or an aluminum silicate containing a Group VIII metal compound and a polyvalent metal halide.\nIt has been found that alkoxy aluminum dichloride catalysts are very useful in the reaction of alkyl esters of acrylic acid with olefins. Alkoxy aluminum chloride catalysts yield excellent results and are much more convenient to handle than conventional ene reaction catalysts such as ethyl aluminum dichloride because they are not pyrophoric."}
{"text": "It is well known in the vehicle industry that certain rotary accessory devices are often used in connection with vehicle motors. Two such accessory devices are alternators and starter motors. Alternators are used in connection with an engine and are typically belt driven by the engine. Alternators have internal components, which when rotated supply electrical power to a vehicle battery. Alternators are typically removable but rigidly mounted via a bracket to the engine block or the chassis of the vehicle. In many cases, where a standard type of alternating mounting arrangement is used, the alternator has “ears” with holes that are mounted onto a post or bolt attached to the vehicle. This permits pivoting of the alternator so that the alternator can be pivoted around the post against the belt tension in order to install and remove belts, and provide a suitable tension when the belt is installed.\nStarter motors are electrical motors, which are typically removable but rigidly mounted to an engine or transmission casing. The starter motor has an electrically driven pinion gear extending from the starter motor that engages a component (typically gears on the flywheel of the engine) in order to be able to rotate the crankshaft of the engine to start it. There is a wide range of attachment mechanisms for attaching the described starter motor.\nConventional alternator and starter tester may test alternators and/or starter motors using a direct current (DC) power source. For example, the DC power source may be used to power the starter motor and bias the alternator during a testing process. The DC power source may determine the capability of an alternator and starter tester. Oftentimes, the alternator and starter tester requires high power in order to achieve desired accuracy when testing the starter motor or the alternator. The high power required leads to expensive and bulky circuit components which in turn lead to higher cost, size and weight for the alternator and starter tester. Additionally, the high power demand required by the prior art systems resulted in greater power consumption and operating costs. Thus, it would be desirable to have an alternator or starter tester with high accuracy with decreased cost, size and weight."}
{"text": "The present invention relates to semiconductor device manufacturing, and more particularly to a method for introducing an impurity gettering layer or region into a top Si-containing layer of a silicon-on-insulator (SOI) wafer.\nIn the semiconductor industry, SOI technology is becoming increasingly important since it permits the formation of high-speed circuits. Moreover, SOI devices offer many more advantages over their bulk counterparts including, for example, higher performance, absence of latch-up, higher packing density and low voltage applications.\nIn a typical SOI wafer, a buried oxide layer is sandwiched between a bottom Si-containing layer, e.g., a Si-containing substrate, and a top Si-containing layer. The buried oxide layer is used to electrically isolate an integrated circuit (IC) fabricated in the top Si-containing layer from the Si-containing substrate. This improves device performance because the buried oxide layer eliminates parasitic capacitance and it allows strong gate to body coupling to improve drive-on current.\nDespite the above advantages, SOI wafers have some problems associated therewith. One problem is that there is an absence of a good gettering layer which removes metallic contamination or impurities in the IC. In a typical bulk semiconductor wafer, the backside of the wafer is typically designed to act as a metal gettering layer. Methods such as providing a heavily doped polysilicon layer on the backside of the wafer or providing a rough backside surface are typically used in bulk semiconductor technology to introduce a gettering layer to the backside of the wafer. In SOI technology, the buried oxide layer blocks metal from getting into the back of the wafer. Thus, an SOI wafer is more sensitive to top silicon metal contamination. Hence, once metal is able to diffuse in the top Si-containing layer of an SOI wafer it has nowhere to go and the metal contamination will damage the device gate dielectric.\nU.S. Pat. No. 5,244,819 to Yue describes a method of providing a frontside gettering layer in an SOI wafer. The gettering layer in Yue is formed by ion implanting species such as Ar, He and Ge into the top Si layer of an SOI substrate. In order to ensure that the gettering layer is not formed in the active device area (i.e., the device channel region) of the top Si layer, a masking layer is used to define the active area. The masking layer is applied to the SOI substrate prior to implantation and a lithographic process is used in defining the masking layer. In accordance with the disclosure of Yue, the non-active device areas (i.e., source/drain regions) of the top Si layer are left exposed and the active device areas are protected by the patterned masking layer.\nU.S. Pat. No. 5,753,560 to Hong, et al. disclose another method of introducing a gettering layer into the top Si layer of an SOI substrate. As was the case in Yue, Hong, et al. use a masking level to ensure that the gettering layer is not introduced into the active device areas.\nAlthough the methods used in Yue and Hong, et al. are capable of providing a gettering layer into the top Si layer of an SOI substrate, the Yue and Hong, et al. processes are not self-aligned, and a costly lithographic process is required in these prior art processes to ensure that the gettering layer is not implanted into the active device regions of the SOI substrate.\nIn view of the above, there is still a need for providing a new and improved method of forming a gettering layer (or region) into non-active device areas of a top Si-containing layer of an SOI substrate that is self-aligned, and which avoids the use of a mask to define the gettering region. A method which eliminates the mask to define the gettering region is beneficial since such a method would also eliminate the costly lithographic process required in prior art methods.\nOne object of the present invention is to provide a method of introducing a gettering region into the non-active device areas of a top Si-containing layer of an SOI substrate.\nAnother object of the present invention is to provide a method of introducing a gettering region into the non-active device areas of a top Si-containing layer of an SOI substrate which is self-aligned.\nA still further object of the present invention is to provide a method of introducing a gettering region into the non-active device areas of a top Si-containing layer of an SOI substrate wherein a mask (i.e., litho level) is not required to define the gettering region.\nThese and other objects and advantages are achieved in the present invention by forming a disposable spacer on each vertical sidewall of a patterned gate stack region prior to forming the gettering region. Because of the spacer, the gettering region will be automatically self-aligned to the patterned gate stack region and it will not impact the active device region, i.e., device channel region. The disposable spacer can then be removed after formation of the gettering region and normal complementary metal oxide semiconductor (CMOS) device processing may be resumed.\nIt noted that in the inventive method the gettering region is introduced into the top Si-containing layer of an SOI substrate after the gate dielectric has been formed thereon. The benefit of such a method is that it does not impact the formation of the gate dielectric. In addition, any metal contamination that is introduced before the gate dielectric is formed could be gettered away using an acid such as HCl immediately before formation of the gate dielectric.\nAnother feature of the present invention is that the gettering region is introduced into the SOI substrate prior to any high temperature processes are performed. Thus, if any metal is present during the high temperature processes no damage to the gate region occurs in the present invention.\nSpecifically, the method of the present invention comprises the steps of:\n(a) forming a disposable spacer on each vertical sidewall of a patterned gate stack region, said patterned gate stack region is formed on a top Si-containing layer of an SOI substrate;\n(b) implanting gettering species into said top Si-containing layer not protected by said disposable spacer and said patterned gate stack region; and\n(c) removing said disposable spacer and annealing the implanted gettering species so as to convert said gettering species into a gettering region."}
{"text": "The invention relates to sealing connectors in gas fastening devices, mounted between the gas cartridge and the solenoid valve for filling the combustion Chambers of these devices in order to ensure sealing of the duct extending through the outlet and filling fittings of the cartridge and the solenoid valve. The term “solenoid valve” is used here for the sake of simplicity, although it could also be a simple mechanical valve and, more generally, any other member for transmitting the gas from the cartridge towards the combustion chamber.\nPossible fastening devices are, e.g. nail drivers for applications with wood.\nThe connectors in question, e.g. such as those described in French patent applications FR 0405605 and FR 0407208 not yet published at the time of filing this application, include a double internal sealing sleeve for receiving the two fittings made of a flexible material, with a generally central transverse annular partition for dimensional adjustment, and an external slide in which the double sleeve is mounted.\nThe slide ensures that the shape or morphology of the flexible double sleeve is maintained and allows the connector to slide in the adapter of the cartridge, by virtue of which, when the cartridge is removed from the device, the sleeve made of flexible material does not remain “stuck” against the adapter and does not hold the fitting of the cartridge in the gas exhaust position.\nOn the one hand, the fitting of the cartridge may pierce the central partition for dimensional adjustment made of flexible material and, on the other hand, this central partition may he deformed to such an extent that it will obstruct the duct for the passage of gas.\nThere is a need for a sealing connector made of one single material sufficiently flexible to ensure sealing between the two fittings, but sufficiently rigid to prevent the central partition from being damaged or deformed and blocking one of the fittings for the passage of gas."}
{"text": "1. Field of the Invention\nThe present invention relates to an arrangement of an oil filter for engine oil and an oil passage in a vehicle engine. More particularly to an arrangement for providing an ease of maintenance.\n2. Description of Background Art\nConventionally, in a vehicle engine provided with exhaust pipes which are curved downwardly from a front wall of a cylinder body, pass in front of and below a crankcase, and extend rearwardly of an engine body, a cartridge type oil filter is available that is attached between a vertically extending portion of one of the exhaust pipes and a front wall of the crankcase so as to be detachable in a vehicle body width direction from a side end side of the engine body. See, for example, Publication of Japanese Patent No. 2705777. Thus, obstructions with respect to the attachment/detachment direction of the oil filter are removed to enhance an ease of maintenance. In addition, front and back portions of the oil filter can be protected by the exhaust pipe and the crankcase.\nThe attachment/detachment of the cartridge type oil filter usually requires a rotation of the oil filter. Accordingly, it is necessary to rotate the oil filter itself by utilizing a tool or an individual's hand.\nHowever, in the conventional vehicle engine described above, the exhaust pipes and the crankcase are arranged in front of and to the rear of the oil filter. Accordingly, it is difficult to position the tool or the individual's hand onto the oil filter. In addition, it is difficult to perform the rotation of the oil filter. Moreover, widening a gap between the exhaust pipes and the crankcase for the purpose of ensuring a working space increases an outer dimension of the engine and lowers a degree of freedom in vehicle body layout. Therefore, it is not preferable."}
{"text": "Demand for information technology devices is constantly growing. Such devices can be equipped with wireless technologies. There is currently a demand for increased performance from these wireless technologies. The demand and general prevalence of the devices has created crowding in the Electro-Magnetic (EM) spectrum. This crowding may be related to EM spectrum partitioning. As more bandwidths are used, there may be fewer available channels. There may be an increasing demand for high-bandwidth usage (e.g., streaming high-definition video to a cellular device). Existing technologies, which allow for such high-bandwidth usage, may reach unsustainable levels of demand and usage. The fact that these technologies, such as Long Term Evolution (LTE) networks, operate on a relatively limited partition of the EM spectrum can magnify these problems. Efficient radio communication in its current state may not be sustainable using current techniques."}
{"text": "Generally, a refrigerator is a home appliance for storing foods in the inside of the refrigerator shielded by a handle thereof at a low temperature, and can store the stored foods in an optimal state by cooling the inside of the storage space using cold air generated by the heat exchange with refrigerant that circulates a freezing cycle.\nThe demand for large and multi-functional refrigerators according to the change and improvement in dietary life has been grown. Taking users' convenience into consideration, products having diverse forms and convenient devices have thereby been released.\nThe inside of the refrigerator may be partitioned by shelves, drawers, baskets, or the like, and the foods may be received in each storage space partitioned thereby. In particular, in the space receiving foods in an airtight state or a state equivalent thereto such as drawers, foods, which are separated and received from other foods, may be stored. Also, foods, which are sensitive to temperature and moisture such as vegetables, are generally stored.\nIn order to store the foods such as vegetables in the storage space such as drawers for a long time, the inside of the storage space may be filled with as little air as possible. To this end, there is provided a configuration to degass air inside the storage space.\nKorean Patent Registration Nos. 0547426 and 0606728 disclose a refrigerator to degass air inside a storage space, wherein a portion of air inside the storage space is thereby exhausted to the outside to reduce the amount of air inside the storage space, so the storage performance of foods can be improved.\nHowever, in Korean Patent Registration No. 0547426, the airtight between a tray and a cover is hardly guaranteed, and the cover must be able to move up and down so as to operate the airtight of the storage space, so an upper portion of the cover must have a space for the operation, thereby causing a problem that the receiving space is lost. Furthermore, efficient close-adhering and degassing can be performed only when the entirety of the cover moves simultaneously, thereby causing a problem that the operation therefor is not easy.\nIn Korean Patent Registration No. 0606728, air is exhausted backward, that is, in a direction that a storage container is drawn in, to have a difficulty in exhausting air smoothly, and the storage container must be compulsorily opened in order that external air enters the inside of the storage container for opening the storage container, thereby causing a problem that use convenience is not good."}
{"text": "1. Field of the Invention\nThe present invention is directed to a wireless Internet service provider, mobile network operator or any such wireless operators for private or public use. Specifically, the invention is adapted to facilitate the recovering of a data communication or data call across different wireless technology platforms or across different private and public wireless networks, provided the user is an authorised user of such different wireless technology platforms or such different wireless-service providers or -network operators such that loss of data on one wireless-platform or -provider and a data reconnection on a same or different wireless-platform or -provider allows recovery of the original data call and thus a continuation of the same original data call.\n2. Technical Background\nCurrent wireless systems and methods do not provide for a way to recover and continue a wireless data communication or wireless data call across different radio technologies. The prior art solutions only apply when a wireless data communication or wireless data call remains within the same radio technology. The most commonly known definitions (mostly defined as standards) of these partial solutions in the prior art are known as different names, such as seamless hand-over, non-seamless hand-over, mid-call hand-over and so forth, however none of these standards and non-standard partial solutions in the prior art resolves the main issue for wireless end-user devices; which is to recover and continue a data call when changing throughout different radio technology platforms available and authorised for data use on such end-users wireless devices.\nAs a matter of illustration to understand the main issue this invention resolves and which has not been fixed by any of the prior arts; say a wireless device is capable of handling data communications or data calls through the following radio data technologies: WiFi (wireless fidelity) through a private- or public-WLAN (wireless local area network), GSM/UMTS (Global System for Mobile Communications/Universal Mobile Telecommunications Service) and that same user's wireless device is an authorised user of such previously mentioned wireless systems that provide radio coverage for such systems (i.e. WiFi and GSM/UMTS) then in the prior art if a wireless device user initiated a data call through WiFi and whilst it's on-going on that 1rd radio technology platform WiFi to a 3rd party and that data call is lost and the wireless device reconnects to a 2nd radio technology platform GSM/UMTS, then such wireless device original data call is terminated permanently in the prior art.\nThe same would happen, namely a data call from a wireless device as described previously would be terminated after the initial data call on a 1st radio technology platform looses data connection or radio coverage and such same data call would not recover and continue when the same wireless device reconnects to a 2nd different to the 1st radio technology platform in any combination; when from WiFi to GSM/UMTS or from GSM/UMTS to WiFi or between any such other wireless technology platforms or also referred herein as radio technology platforms.\nRecently, fixed line and mobile network operators and service providers in particular have begun to offer on-line (internet) services allowing individuals to make wireless text-, voice- and video-data calls over a variety of standards or proprietary protocols competing head-on with the traditional analogue and digital data calls from traditional network operators, yet none provide a solution that would allow a wireless device end-user to restore an original data call on one radio technology to continue that same data call when the wireless device reconnects on another radio technology, provided the time between loosing data connection from the 1st radio technology and re-connecting to the same or to a 2nd radio technology is within a reasonable timeframe acceptable to end-users (i.e. in the region of seconds or tenths of seconds not minutes).\nAttempts have been made to ease the way to wireless end-users always focused on some form of hand-over defined in standards and some non standard; however none of the prior art allows any reliable solution across different radio technologies available in end-users wireless devices to which they are authorised users to, other then the partial solution described above, which in itself does not provide a solution to the issue described herein before."}
{"text": "The reason communications revolution has increased the amount of data available to the public. The public's hunger for more and more information has been whetted by this excess of information.\nCurrent systems for information dissemination, especially for disseminating to the members of the public who take public means of transportation, are currently inadequate. These systems are solely visual and static. Sandwich boards, billboards, and other very old static forms of advertising are used for such ends. However, as noted above, these systems are inadequate. Such static forms are purely visual and cannot be readily changed in response to changed conditions. Unfortunately, the ubiquity of televisions and computers have had a deleterious effect on the attention span of the average person. As such, current static systems are unable to keep people's attention.\nAnother major drawback of current systems is the lack of an audio component. Again, the proliferation of multi-media devices and the ubiquity of multi-media in today's society has the public clamouring for the same presentation type in almost all of their information input. While there are some information dissemination systems that integrate an audio portion (such as publicly accessible television sets), these can be intrusive and annoying to those who do not wish to take advantage of the system.\nThere is therefore a need for an information dissemination system that integrates both video and audio components while simultaneously avoiding becoming an annoyance to those not wishing to access the information being disseminated. Such an information dissemination system should also be dynamic in that modifying the information content should be relatively simple. Furthermore, the system should be able to keep end users (i.e. the public, especially those who are travelling in public transportation vehicles) occupied and should be able to hold their attention. Finally, it is also very desirable if such a system would be cost effective and financially self sustaining. It is therefore an object of the present invention to mitigate if not overcome the deficiencies of the prior art."}
{"text": "The referenced U.S. Pat. No. 4,425,888, Engel, et al., assigned to the asignee of the present application describes a speed control or governing system for an automotive-type internal combustion engine, especially a self-ignited internal combustion engine, that is, for example of the diesel engine type. In accordance with the disclosure of the reference, the command value of speed is raised by a proportional-integrating acting controller when the actual speed has reached a predetermined relationship with respect to the command speed. In such a speed control system, commanded speed is changed when the difference between actual speed and commanded speed reaches a predetermined level or value. The difference between actual speed and commanded speed is maintained constant.\nAs the speed drops, particularly when the speed drops below a lower threshold level, a different control function becomes effective. The drop in commanded speed is then matched, at least roughly, to the actual speed drop, as determined by operating characteristics of the engine when it has normal operating temperature.\nIf the engine is cold, the actual speed drop may be more rapid than the speed drop when the engine is warm. Consequently, it is possible that the controller will not follow rapidly enough the actual changes in speed of the engine and a predetermined minimum idle speed may actually be passed, by the engine when it is still cold and the engine will stall. The control system, as described, while functioning well with an engine which is warm may not prevent stalling of the engine due to the excessively rapid speed drop when the engine is still cold."}
{"text": "The present invention relates to controls for hydrostatic transmissions and more specifically relates to a mechanism for neutralizing a hydrostatic traction drive transmission in response to actuation of a traction drive brake or clutch.\nIt is known to interlock manually operable traction drive brake or clutch controls with those for controlling the direction and speed of a hydrostatic traction drive transmission so that the transmission is placed in a neutral condition, and a speed and direction control lever is placed in its neutral position in response to movement of levers for effecting brake engagement or clutch disengagement. One example of such interlocked controls is disclosed in U.S. Pat. No. 3,511,105 granted to Matter on May 12, 1970.\nThe known interlocked controls suffer from one or more of the drawbacks of being relatively complex, of requiring specially configured cam slots for effecting simultaneous neutralizing of the transmission control arm and control lever and of requiring frequent adjustment of control linkage elements to ensure such simultaneous neutralizing."}
{"text": "The discovery, development and commercialization of the 2-(2-imidazolin-2-yl)pyridine and quinoline compounds as herbicidal agents has given new meaning to the term \"weed control\"; for within this series of compounds it has been found that some are broad-spectrum or total vegetation herbicides with activity in both herbaceous and woody plants. Others are highly selective weed control agents useful as weed control agents in the presence of crops.\nSeveral processes for the preparation of the herbicidal 2-(2-imidazolin-2-yl)pyridines involve the preparation of 2,3-pyridinedicarboxylic acids from 2,3-pyridinedicarboxylic acid diesters, but methods for the preparation of 2,3-pyridinedicarboxylic acid diesters produce impure products. Current methods of production of 2.3-pyridinedicarboxylic acid diesters by various condensation routes provide very impure products (20% to 50% purity). Arduous or time-consuming purification methods are required to provide purified 2.3-pyridinedicarboxylic acid diesters. Distillation has been used to improve the purity initially to as high as about 65% to 70%. However, the distillation is done at high temperatures and at high vacuum and the distillation product is further purified by a second distillation to obtain purified 2,3-pyridinedicarboxylic acid diesters.\nSince the purity of the 2,3-pyridinedicarboxylic acid diesters impact directly on the purity of the 2,3-pyridinedicarboxylic acids, a process that would improve the purity of the 2.3-pyridinedicarboxylic acid diesters without requiring the use of time-consuming distillation procedures would provide a great improvement in the processes used to prepare the herbicidal 2-(2-imidazolin-2-yl)pyridine and quinoline compounds.\nIt is an object of the present invention to provide a process for the purification of 2,3-pyridine and quinolinedicarboxylic acid diester compounds which avoids the use of the arduous or time-consuming purification methods of the prior art."}
{"text": "a. The Field of the Invention\nThis invention relates to the field of network communications. In particular, the invention relates to a network where devices communicate using sinusoidal waves.\nb. Background Information\nEthernet is an industry standard (e.g., IEEE 802.3 specification) method of communicating between various devices in a local area network. For example, a computer includes a network interface card (NIC) that formats ethernet data for transmission onto a network cable. The network cable carries the ethernet formatted packets out to the rest of the network. The data signal is generated to comply with a particular specification for that type of ethernet communications. For example, the ethernet data signal might be generated to comply with the ANSI/IEEE standard 802.3 ethernet voltage template. This voltage template applies to five and ten MHz ethernet data communications. Complying with the voltage template ensures that the NIC will not damage other devices connected to the network cable, and ensures that other devices will be able to properly receive and decode the signals from the NIC.\nAs part of the FCC's (and other national regulatory bodies) electromagnetic interference regulations, the NIC must not emit an amount of electromagnetic radiation above a preset limit. The FCC (or for Europe, the CISPR-B) specification, for example, limits the radiation from the ethernet NIC. Importantly, any multiples of the fundamental frequency are to be less than twenty-seven DB lower than the fundamental frequency, per the 802.3 specifications. For example, when driving an all-ones Manchester encoded signal, any harmonic measured shall be at least twenty-seven dB below the fundamental. One problem with generating the ethernet signal to comply with the ethernet voltage template is that high frequency noise may also be generated. This high frequency noise typically appears as an unacceptably high electromagnetic radiation signal at the connection between the NIC and the network cable. These high frequency signals can cause the network interface card to not comply with the FCC electromagnetic interference requirements.\nOne prior art NIC is shown in FIG. 1. This NIC is available from 3COM Corporation of Santa Clara, Calif. The NIC 100 is for generating the transmit signal 130 which corresponds to an ethernet transmit signal. The transmit signal is generated to support Manchester encoding of the ethernet data. The transmit signal 130 is approximately a square wave signal. The prior art system of FIG. 1 uses a square wave in combination with a pre-emphasis wave. Each of these waves are ten MHz and staggered five MHz but the combination provides a partially shaped wave form as shown in the transmit signal 130. The partially shaped wave form helps to meet the voltage template of the transmit signal 130.\nThis paragraph describes the operation of the NIC 100. The NIC 100 includes an ethernet controller 101, some filtering circuits, and a transformer. The ethernet controller 101 generates ethernet data signals and includes a number of digital output signals. These output signals include the twisted pair pre-emphasis minus 102, the twisted pair data plus 104, the twisted pair pre-emphasis plus 106, and the twisted pair data minus 108. The twisted pair data plus 104 and the twisted pair data minus 108 output signals that correspond to the positive portion and the negative portion of the transmit signal 130. The twisted pair pre-emphasis minus 102 and the twisted pair pre-emphasis plus 106 help shape the transmit signal 130 to better comply with the voltage template of the IEEE 802.3 specification. The circuitry between the outputs of the ethernet controller 101 and the transformer 120 helps combine the outputs into the one transmit signal 130. The transformer 120 includes a seven pole filter which shapes the transmit signal 130 and reduces the high frequency components in the transmit signal 130.\nThe NIC 100 operates quite well for many applications. However, the external circuitry on the NIC 100 adds to the costs of manufacturing the NIC 100. Therefore, it is desirable to generate the transmit signal with few external components. Additionally, the electromagnetic interference should be kept to a minimum."}
{"text": "It is well know that edible oils deteriorate during cooking, particularly when they are repeatedly heated to high temperatures. These oils are typically heated to temperatures of the order of 180° C. to fry food. A multitude of chemical reactions occur at these temperatures, such as polymerisation, thermo-oxidation, etc., which significantly alter the quality of the oil. The quantity of some products of these reactions must not exceed a threshold imposed by legislation, since the oil is deemed unfit for consumption beyond the threshold. It is thus important to be able to detect the threshold in a reliable manner, so that the oil is replaced as soon as it becomes necessary. For a long time, it was left to cooks to judge, after a visual and/or olfactory inspection, whether the oil was still fit for consumption. Of course, that method is entirely subjective and is consequently unreliable.\nEP Patent No, 1 588 158, which corresponds to U.S. Patent Application Publication No. US 2006/0288877 A1, discloses a device for the capacitive measurement of the quality and/or deterioration of a cooking oil to overcome these drawbacks. The content of EP Patent No, 1 588 158 and corresponding U.S. Patent Application Publication No. US 2006/0288877 A1 is incorporated herein by reference. In this device, the capacitive sensor is directly arranged in the vat of the cooking apparatus, with the sensor encapsulated in a perforated protective case, secured in a submerged area of the vat.\nAlthough the device disclosed in that Patent Application operates satisfactorily, performing a capacitance measurement inside a deep fat fryer remains a highly delicate operation. Indeed, a variation of a few picofarads between the new oil and used oil, greatly influenced by temperature, water and impurities, is not easy to detect. Added to this is the fact that the device has to operate in a very harsh environment in which the capacitive sensor, and, possibly, the temperature sensor associated therewith, are subjected to temperatures higher than 200° C., and to shocks when careless operators strike the sensor with the baskets holding food for frying.\nProtection of the capacitive sensor, and possibly, the temperature sensor associated therewith, thus constitutes an extremely important problem that needs to be addressed to avoid deterioration in measurement accuracy and/or reducing the life time of the measuring device, since the measurements are directly dependent thereon."}
{"text": "This invention relates to image processing and, more particularly, to the characterization of images.\nImage processing refers to the analysis and manipulation of images using a computer. The discipline encompasses a wide array of techniques, including pattern recognition and other image analysis, compression and encoding, such as image transmission, and image construction, to name but a few examples. Image processing is used in diverse fields such as astronomy, medicine, military operations, communications, geology, meteorology and so on.\nAlthough visual in nature, images may be characterized so that a computer or other processor-based system may also xe2x80x9cseexe2x80x9d the image or, for example, distinguish the image from other images. Identifying the moment, entropy, center of mass, orientation, or other histogram-based features is one approach to image characterization. Structural and syntactic features of the image may serve as characteristic parameters. Geometric indicia, such as perimeter, area, eccentricity, or Euler number, may numerically identify the image.\nBecause images are often perturbed, such as during transmission, the characteristic parameters of the image are ideally invariant following transformations. Image transformations may include translation, rotation, scaling, shearing, compression, inclusion of noise, and so on.\nParticularly in the emerging domain of Internet technology, image searching may be an important tool. Different characteristics of the images, used for image searching, may have widely varying results, however. This type of image searching is known as content-based image retrieval.\nOne technique for retrieving images from a database of images is to associate each image with a certain number of features. Using feature vectors, each image, represented by its feature values, is mapped to a point in an n-dimensional feature space, where there are n features identified. A similarity query is then used to retrieve the images lying in a neighborhood close to the query image.\nIn executing these similarity queries, the precision of obtaining the image may be offset by the size of the image database and the time to calculate the various features involved in the similarity query.\nThus, there is a continuing need for optimizing image retrieval methods."}
{"text": "Communication between electronic devices is becoming ubiquitous as data communication networks expand into the very fabric of our society. Once limited to the laboratory, data communication networks now connect nearly every imaginable electronic device, including large mainframe computers, mini or microcomputers, personal or handheld computers, personal digital assistants, and electrical controllers such as light switches and thermostats. Even devices that use simple microprocessors with small amounts of memory can operate over limited-bandwidth networks. Using a proper interface, many of these electronic devices can connect to larger communications networks over standard phone lines or cable television lines. Wireless networks allow portable electronic devices, such as personal digital assistants, cellular telephones, or interactive pagers, to connect to more traditional, wired networks. Thus, the ease with which an average user manipulates his or her networked, electronic environment bears directly on the frequency of use and ultimate satisfaction derived from advances in technology.\nFor example, call forwarding has been a common feature of communication networks for a number of years. And yet, calls are frequently not forwarded for a variety of reasons, including user location uncertainty, security concerns, or inconvenient command requirements.\nOne device known in the prior art is U.S. Pat. No. 5,197,092 to Bamburak. Bamburak discloses a call forwarding notification system using a receiving station (or “holster”), including a standard phone connected to a landline system (i.e., the PSTN), and a personal communicator (i.e., a cellular telephone) connected to a wireless system. When the cell phone is placed in the receiving station, the receiving station's phone number is sent to the wireless network central exchange, using the call forwarding update number so that all future calls to the cell phone are forwarded to the receiving station over the landline system (i.e., the cell phone's call forwarding feature is turned on). Bamburak also discloses that the receiving station's phone number may be sent to the wireless network central exchange via either the landline system or the wireless system. However, the receiving station and cell phone must be able to communicate with one another via a two-way interface. In Bamburak's “best mode”, the cell phone's call forwarding update number and system identification (SID) number are stored in the cell phone and communicated to the receiving station, while in Bamburak's “variation”, the receiving station's phone number is stored in the receiving station and communicated to the cell phone.\nBamburak's call forwarding notification system, as well as other systems known in the prior art, fails to satisfy the needs of the average user when confronted with task of controlling his networked, electronic environment. For example, Bamburak fails to discuss whether call forwarding is turned off when the cell phone is removed from the receiving station. Bamburak also fails to teach or suggest that other telecommunications features, such as voice messaging, may be automatically controlled or even that other network devices, such as home automation devices or security systems, can be conveniently controlled. Furthermore, Bamburak fails to disclose whether an active phone call can be transferred from the cell phone to the receiving station when the cell phone is placed in the receiving station, nor does he address whether this transfer can be accomplished without an incoming ring signal. And Bamburak discusses call forwarding in the context of a single phone number (i.e., the cell phone) and fails to consider the advantages of conveniently manipulating all of the phone numbers typically associated with an average user (e.g., cell phone, home phone, multiple work phones, etc.).\nThus, there is a need in the art for a simple and intuitive method for notifying all the appropriate network communications devices that a user has arrived at a particular place, and for those devices to automatically reconfigure the user's communication infrastructure in an appropriate way."}
{"text": "1. Field of the Invention\nThe present invention relates to a brushless motor.\n2. Description of Related Art\nJapanese Unexamined Patent Publication No. 2008-86064A (corresponding to US 2008/0073995A) recites a brushless motor known as a 10-pole/12-slot brushless motor. In this brushless motor, teeth of a stator core are arranged one after another at 30 degree pitches, i.e., intervals (360 degrees/12=30 degrees) in the circumferential direction of the brushless motor, and magnetic poles of rotor magnets are arranged one after another at 36 degree pitches, i.e., intervals (360 degrees/10=36 degrees) in the circumferential direction.\nIn the above-described brushless motor, a difference (specifically, a difference of 6 degrees measured as the central angle) exists between the pitch of the teeth measured in the circumferential direction of the brushless motor and the pitch of the magnetic poles of the rotor magnets measured in the circumferential direction of the brushless motor. This difference poses the following disadvantage.\nGenerally, in order to implement a relatively high effective magnetic flux amount and a relatively high motor efficiency in the motor, desirably, a phase difference Δθ (see FIG. 5) between the phase of the electric current I flowing through the stator coil wound around the corresponding tooth and the phase of the magnetic flux φ applied from the rotor magnet to the stator coil should be 90 degrees, as is well known according to Fleming's left-hand rule.\nAs discussed above, in the case of the brushless motor recited in Japanese Unexamined Patent Publication No. 2008-86064A, the difference (the difference of 6 degrees measured as the central angle) exists between the pitch of the teeth and the pitch of the magnetic poles of the rotor magnets. Therefore, the phase difference between the phase of the electric current flowing through the stator coil wound around the corresponding tooth and the phase of the magnetic flux applied from the rotor magnet to the stator coil becomes 75 degrees, which is smaller than 90 degrees. Thereby, the effective magnetic flux amount may be disadvantageously reduced to reduce the motor efficiency."}
{"text": "This invention relates to cubic boron nitride coated material useful as a tool, such as a cutting tool, a wear resistant tool and the like and an electronic part such as a semiconductor substrate and the like, and the producing method thereof.\nRoughly classifying, boron nitride (BN) is present in two types of forms, one being low-density boron nitride and the other high-density boron nitride. Of these, as a representative example of high-density boron nitride, there is cubic boron nitride which is synthesized under special conditions such as high-pressure and high-temperature.\nCubic boron nitride has high hardness, high heat conductivity and high electric insulating property, second to diamond, and moreover, enjoys chemical stability, oxidation resistance, heat resistance and thermal shock resistance all better than diamond. While diamond has high affinity with Iron-group metals, cubic boron nitride has low affinity with them. Therefore, cubic boron nitride is paid attention, for example, as a material for tools used for cutting or grinding Iron-group metal materials. In spite of such excellent properties as mentioned above, cubic boron nitride is brittle and has low sinterability and limitations are hence imposed on its shape and applicaiton.\nThus, it has been attempted to form cubic boron nitride as a coating layer on a surface of a substrate in order to solve the above-mentioned limitations to the shape and application.\nFormation of coating layers of boron nitride can generally be carried out, roughly speaking, by chemical vapor deposition method (CVD method), physical depostion method (PVD method) and plasma-assisted CVD method (PCVD method). Of these, CVD method is effected by using, as illustrative, reactant gases comprising a boride such as a boron halogenide or diborane, hydrogen, and ammonia or hydrazine. PVD method includes, for example, the ion beam deposition method, ion injection method, sputtering method and ion plating method as well as combinations of the ion injection method with other PVD methods. In plasma-assisted CVD method, a coating layer of boron nitride is synthesized in a vapor phase in the stream of plasma. As a prior art publication disclosing tools obtained by forming coating layers actually consisting of cubic boron nitride on substrates, Japanese Provisional patent publication No. 95881/1982 has been known.\nAmong conventional methods for forming coating layers consisting of boron nitride, CVD method is accompanied by a problem that coating layers comprising of hexagonal boron nitride or amorphous boron nitride can only be formed because it is a mere thermal vapor-phase synthesis. On the other hand, when PVD method or plasma-assisted CVD method is relied upon, a coating layer composed of boron nitride having a relatively high hardness is formed. Cubic boron nitiride has been considered to be the component of this coating layer. Coating layers formed by PVD or plasma-assisted CVD method are, however, accompanied by such drawbacks that they contain cubic boron nitride at low contents and due to the overall low hardness of the coating layers and the weak bonding between the coating layers and their corresponding substrates, the resultant coated layers cannot be used for actual applications.\nIn Japanese Provisional patent publication No. 95881/1982, which discloses tools produced in accordance with these conventional methods, coating layers made of cubic boron nitride are formed directly on surfaces of sintered silicon nitride bodies as substrates. They are produced by conventional CVD, PVD or plasma-assisted CVD method or alternatively by coating the surfaces with hexagonal boron nitride and then heat-treating the coating layers of hexagonal boron nitride to form coating layers consisting of cubic boron nitride. Therefore, the above-mentioned coated tools are accompanied by such drawbacks that even if coating layers of cubic boron nitride are formed on the surfaces of the substrates, the contents of cubic boron nitride in coating layers are very low or otherwise the adhesion between substrates and their corresponding coating layers are poor."}
{"text": "In a common design of pneumatic conveying systems the materials being conveyed and the conveying gas together pass through a materials handling fan located in a pneumatic conveying line between the material pick-up or feed point and the material discharge point.\nFans which only handle gases normally have impeller blades attached to a circular front and rear shroudplate. In such a fan gas passes from the inlet opening at the center of the impeller radially outwards in the constricted paths defined by the blades and the front and rear shroudplates to the periphery of the impeller where it leaves the fan casing via an exit opening.\nIn designing a fan handling only gases or in a materials handling fan the principal design objective is to get a specific throughput. That throughput is a function of pressure difference across the fan and volume of gas passing through the fan. These parameters, in turn, are a function of fan blade tip speed and blade width, respectively. Blade tip speed is, of course, a function of the blades' radius and rotational speed. Thus, for a given rotational speed the pressure difference across the fan can be increased by increasing the fan blade length or impeller radius. Similarly, to increase volume through the fan the width of the fan blade can be increased. Of course, as these dimensions increase the stresses on the impeller increase dramatically for the reasons described in more detail below.\nThe double shroud design of the previously described gas handling fan is not suitable for use as a materials handling fan because the conveyed material cannot readily pass through the constricted path defined by the blades and the front and rear shroudplates.\nThe problems associated with pneumatically conveying materials through a fan with a double shroudplate impeller have led to the design of materials handling fans with impellers having radial blades shroudplate. These open impellers are not as efficient as the double shrouded gas handling impellers because of the gas turbulence that takes places between the stationary fan casing and the unshrouded side or sides of the radial impeller blades. The materials handling fan impellers with one shroudplate are substantially more efficient than the impellers with no shroudplate and the single shroud design is therefore most favored and is almost always used in the larger diameter, higher rotational speed fans used in systems requiring a substantial generation of operating gas pressure.\nThis invention deals only with the more efficient materials handling fan impeller using a single shroudplate or backplate.\nThe single shroudplate radially bladed materials handling impeller design is subject to unusually high stresses near the hub on the open or materials feeding side of the impeller because of its asymmetric design. During rotation the centrifugal force acting on the impeller blades and the shroudplate causes stresses and strains in both these components. At the same rotational speed the strain on the impeller blade alone will be substantially greater than the strain on the shroudplate of equal diameter because of the constraint of the circumferential ligaments in the circular shroudplate. As a specific example the radius of a 20\" radius .times. 3/8 thick steel disk used as a shroudplate rotated about its central axis at 3600 rpm will increase approximately 0.005 inches, whereas the radius of a 20\" long .times. 3/8 thick radial blade rotated about one end at 3600 rpm will increase approximately .009 inches.\nWhen a materials handling fan impeller, made up of a multiplicity of radial blades welded or otherwise affixed to a single circular shroudplate, is rotated about its central axis, the greater strain or radial growth of the radial blades as compared to the strain or radial growth of the shroudplate will cause the impeller to distort so that the impeller, when viewed from the radial blade side, will be convex with the highest stresses being present on the non-shroud side of the impeller closest to the impeller hub.\nThese high stress levels, compared to the stress levels in a similar sized double shroud impeller operated at the same speed, have dramatically limited the diameter and operating speed of commercial materials handling fan rotors with a single shroudplate and hence their performance capabilities.\nSophisticated designs of materials handling fan impellers of large diameters for high rotational speeds which incorporated tapered blades and back plates, large specially shaped hubs, etc., have been developed which have the materials of construction strategically located so that the large centrifugal forces present during rotation generate acceptably low and safe levels of stress. These impellers require either extensive and expensive machining in order to distribute the materials satisfactorily or, alternatively, large and expensive castings.\nIt would therefore be desirable if an impeller could be constructed which is basically adaptable for use with fans of varying diameters. Absent careful design of such impellers, however, high failure rates for fans of larger diameters are likely because of the inability of the impeller to withstand stress at the impeller hub.\nA background patent of general interest in this area is U.S. Pat. No. 4,285,635 which discloses a centrifugal-blower impeller -- not a materials handling fan -- having a center part and a cicumferential outer part of lesser thickness than the center part. The outer part consists of ring segments welded to each other and the periphery of the center part to create compressive forces in the center part."}
{"text": "The present invention relates to an apparatus for detecting the sudden evolution of an acoustic signal used for acoustic monitoring and in particular for leak detection. It is more especially used in inspection installations for fast neutron reactors, where it makes it possible to detect leaks in the heat exchange wall of a sodium-heated steam generator. However, it can be used whenever it is desired to detect a leak in a random member.\nIn a fast neutron nuclear reactor, the steam generator generally comprises steel tubes, within which water and steam circulate under high pressure. Liquid sodium carrying the heat given off by the reactor core circulates outside the tubes. The tubes form the heat exchange walls between the water and the sodium. The water evaporates within the tubes and the thus formed steam is used for driving turbines.\nDue to the high chemical affinity between water and sodium, any sealing defect in a tube must be rapidly detected. Thus, any leak of water into the sodium can on the one hand evolve suddenly (leaking at a rate of a few to one hundred grams per second) producing a high pressure field and, on the other hand, can seriously damage the generator through the perforation of adjacent tubes under the action of an erosion--corrosion phenomenon.\nAt present, two sealing defect detection means are used, constituted by means for detecting hydrogen (in the sodium and in the argon) and safety diaphragms.\nA hydrogen detection apparatus continuously measures the hydrogen concentration in the sodium at the outlet from equipment, or in the argon of the covering gases. The sensitivity is excellent, because such an apparatus is able to detect leaks having a flow rate of less than 1 g/sec. However, the response time is long and, in certain cases, may be excessive compared with the breakdown or perforation time of the tubes adjacent to the leak.\nSafety diaphragms serve to prevent excessive overpressures within steam generators. They can be fractured by pressure waves produced in the apparatus by the sudden appearance of a leak and by the oscillation of the sodium mass between the argon volumes present in the secondary loop.\nMore recently a new leak detection procedure has been discovered, which uses acoustic means operating on the following principle. Acoustic sensors are placed on the steam generator to be controlled. Any leak of pressurized steam into the liquid sodium causes the generator walls to vibrate and this is superimposed on the vibration emanating from the different noise sources, i.e. outflow of sodium, outflow of steam, movements of various mechanical parts, equipment vibrations, relaxation of metal stresses, etc. All these vibrations are transmitted through the sodium and the metal parts up to the steam generator walls. The sensors supply a complex electrical signal which reflects all these waves. Its spectral range covers the audible band. The sound wave due to the leak gives rise to a characteristic signal, which is detected by an electronic system connected to the sensors.\nSuch a procedure is described, for example, in the reports of the ANS Conference, held at Richmond (WA), U.S.A., on Apr. 20-24, 1980, and particularly that of D. A. GREENE, F. F. AHLGREN and D. MENEELY entitled \"Acoustic Leak Detection System for Sodium Heated Steam Generators\".\nAlthough this procedure has advantages, it is difficult to perform for the following reasons. The background noise detected by the acoustic sensors has an intensity depending on a large number of parameters and particularly the power of the reactor. Thus, if the power passes, e.g., from 20 to 80% of its maximum value, the background noise can be increased by 20 dB, which is considerable and largely masks an increase in the acoustic signal resulting from a leak. Thus, considerable importance is attached to the processing of the acoustic signal, if it is wished to reliably detect a leak and deliberately trigger an alarm."}
{"text": "In mobile communication systems, the nature of mobile devices creates situations where as a wireless transmit/receive unit (WTRU) moves through an area, the WTRU may encounter multiple gateways through which it could connect. These multiple gateways may use different communication architectures to connect devices and to Internet protocol (IP) access.\nSome gateways may be under the control of a service to which the WTRU user subscribes. Other gateways may belong to systems that are not aware of the user but may still allow the user's WTRU to establish a connection.\nAs the user moves throughout the coverage areas of these gateways, it may be better for the WTRU to switch from a current communication architecture to another architecture, such as a trusted network, or a network that may be able to take advantage of more sophisticated capabilities of the WTRU. When this occurs, a handover takes place to move the access from the original network architecture to the newly detected architecture. Likewise, when a trusted network architecture's signal becomes weak, the WTRU may decide to handover to a different network architecture.\nIn prior handoff methodologies, the WTRU could establish connectivity with the new network architecture and abruptly sever its connection with the original network. It would be beneficial if the original network connection could be terminated in an orderly fashion when handover between network architectures occurs.\nThe network architecture showing the relationship between a third generation partnership project (3GPP) architecture and a non-3GPP architecture is depicted in FIG. 1. A network architecture 100 includes a 3GPP and non-3GPP system architecture, divided by a dashed line 101. Above the line 101 is 3GPP compliant architecture and below the dashed line 101 connections is a non-3GPP architecture. WTRU 103 may gain access to the 3GPP architecture through connection S2a 105, S2b 107, or S2c 109, depending on the architecture to which the WTRU 103 is connected and the relation of that architecture to the 3GPP network. If the non-3GPP architecture is a trusted non-3GPP IP connection 111, the connection to the packet data network gateway 113 is made directly through S2a 105. S2a 105 provides the user plane with related control and mobility support between trusted non 3GPP IP access and the packet data network (PDN) Gateway (GW) 113.\nWhen the non-3GPP architecture is untrusted 115, the connection is made through an evolved Packet Data Gateway (ePDG) 117. The connection between the ePDG 117 and the PDN GW 113 is made through an S2b 109 connection. S2b 109 provides the user plane with related control and mobility support between evolved packet data gateway (ePDG 117) and the PDN GW 113.\nConnection between the WTRU 103 and the PDN GW 113, while the WTRU 103 is connected to either a trusted or untrusted non-3GPP or 3GPP, may be provided through S2c 107. S2c 107 provides the user plane with related control and mobility support between a wireless transmit/receive unit (WTRU) 103 and the PDN GW 113. This reference point is implemented over trusted and/or untrusted non-3GPP access and/or 3GPP access.\nAn S5 connection 119 exists between the PDN GW 113 and a serving gateway 121 in the 3GPP system. S5 119 provides user plane tunneling and tunnel management between Serving GW and PDN GW. It is used for Serving GW relocation due to mobility and in case the Serving GW needs to connect to a non-collocated PDN GW for the required PDN connectivity.\nAn S6a interface 123 is defined between mobility management entity (MME) 125 and home subscriber server (HSS) 127 for authentication and authorization.\nThe point defined by S6c 129 is the reference point between PDN GW 113 in a home public land mobile network (HPLMN) and a 3GPP authentication, authorization and accounting (AAA) server 131 for mobility related authentication if needed. This reference point may also be used to retrieve and request storage of mobility parameters.\nReference point S6d (not shown) is between Serving Gateway in a visited public land mobile network (VPLMN) and a 3GPP AAA Proxy for mobility related authentication if needed. This reference point may also be used to retrieve and request storage of mobility parameters.\nConnection S7 133 provides transfer of quality of service (QoS) policy and charging rules from policy and charging rules function (PCRF) 135 to policy and charging enforcement point (PCEF) (not shown).\nSGi 117 is the reference point between the PDN Gateway 113 and the packet data network 139. The packet data network 139 may be an operator external public or private packet data network or an intra operator packet data network, e.g. for provision of IP multimedia subsystem (IMS) services. This reference point supports any 3GPP and non-3GPP access systems.\nWa* 141 connects the untrusted non-3GPP IP access with the 3GPP AAA server/proxy 131 and transports access authentication, authorization and charging-related information in a secure manner.\nTa* 143 connects the trusted non-3GPP IP Access 111 with the 3GPP AAA server/proxy 131 and transports access authentication, authorization, mobility parameters and charging-related information in a secure manner.\nThe Wm* 145 reference point is located between 3GPP AAA Server/Proxy 131 and ePDG 117 and is used for AAA signaling, (transport of mobility parameters, tunnel authentication and authorization data).\nWn* 147 is the reference point between the untrusted Non-3GPP IP Access 115 and the ePDG 117. Traffic on this interface for an initiated tunnel has to be forced towards ePDG 117.\nWx* 149 is the reference point is located between 3GPP AAA Server 131 and HSS 127 and is used for transport of authentication data.\nA diagram illustrating the handover process between untrusted non-3GPP IP access using PMIPv6 to an E-UTRAN 3GPP network is depicted in FIG. 2. The WTRU 201 is initially connected to the untrusted non-3GPP access. There is an IPsec tunnel 203 between the WTRU 201 and the ePDG 205 and a PMIPv6 tunnel 207 between the ePDG and the PDN GW 209. As the WTRU 201 moves, it may switch from the untrusted non-3GPP IP access to a 3GPP architecture such as E-UTRAN. The WTRU 201 attaches to the E-UTRAN network 211. Next, the WTRU 201 performs access authorization with the mobility management entity (MME) 213. The MME 213 contacts the home subscriber server (HSS) 215 for authentication of the WTRU 201. As part of the authentication procedure, the packet data network gateway (PDN GW) 209 that will be used is conveyed to the MME 213. The MME 213 performs a location update procedure and subscriber data retrieval 214 from the HSS 215. After the MME 213 authenticates the WTRU 201, it sends a create default bearer request message 217 to the serving gateway (GW) 219. In the message 217, the MME 213 includes the WTRU's 201 identifier (NAI) and the PDN GW 209 that will be used. The serving GW 219 sends a proxy binding update (BU) message 221 to the PDN GW 209 to setup a default bearer for the WTRU 201. The proxy BU 221 includes the WTRU's 201 identifier and a request for an IP address from the PDN GW 209. The PDN GW 209 processes the proxy BU message 221 from the serving GW 219, updates the binding cache entry for the WTRU 201 and responds with a proxy binding acknowledgement 223. In the proxy binding acknowledgement (Ack) 223, the PDN GW 209 replies with the same IP address or prefix that was assigned to the WTRU 201 earlier. At that point a PMIPv6 tunnel 225 exists between the PDN GW 209 and the serving GW 219. The serving GW 219 responds to the MME 213 with the create default bearer response message 217. In this message 217, the serving GW 219 includes the IP address of the WTRU 201. An S1_U default bearer establishment procedure 227 is performed. This procedure includes a radio bearer setup 228. At the end of the handover procedure 200, there is a default bearer for the WTRU 201 that consists of E-UTRAN radio bearer 227, 51 bearer between the e Node B and the serving GW 219 and a PMIPv6 tunnel 229 between the serving GW 219 and the PDN GW 209.\nFIGS. 3A and 3B are an illustration of a handover procedure over interface S2c of FIG. 1 from non-3GPP IP Access to 3GPP access. The session starts in untrusted non-3GPP access 301 using DSMIPv6 over the S2c interface. The session hands over to 3GPP access, for example, E-UTRAN 303. The WTRU 305 uses an untrusted non-3GPP access system 301. It has an IPsec/IKEv2 session with the ePDG 307 and a DSMIPv6 session with the PDN GW 309. The WTRU 305 discovers the 3GPP access system 303 and decides to handover from the currently used trusted non-3GPP 301 access system to the discovered 3GPP access system 302. The WTRU 305 sends an attach request 313 which is routed by 3GPP access system 303 to an MME 311 instance in the evolved packet core (EPC) which is not shown. The MME 311 contacts the HSS/3GPP authorization and authentication (AAA) 315 and authenticates the WTRU 305. As part of the authentication procedure, the IP address of the PDN GW 309 that will be used for 3GPP access is conveyed to the MME 311. After successful authentication, the MME 311 performs a location update procedure with HSS 315. The MME 311 selects a serving GW 317 and sends a create default bearer request message 319 to the selected PDN GW 309. The serving GW 317 may initiate the PMIPv6 registration procedure towards the PDN GW 309 by sending a Proxy Binding Update 321, such as when using Internet Engineering Task Force (IETF) based S5 interface between the PDN GW 309 and the serving GW 317. If GPRS Tunneling Protocol (GTP) is used for S5, the serving GW 307 send a create bearer request message 309 to the PDN GW 309. In IETF based S5, the PDN GW 309 responds with a proxy binding acknowledgement 327 and updates its mobility binding which effectively switches the DSMIPv6 tunnel from the non-3GPP access network to the PMIPv6 tunnel to the serving GW 307. In the proxy binding acknowledgement 327, the PDN GW 309 includes the home IP address or prefix that was assigned to the WTRU 305 earlier. For GTP-based S5, the PDN GW 309 responds with a create bearer response message 329 to the serving GW 307. The create bearer response 329 contains the home IP address or prefix that was assigned to the WTRU 305 earlier. The serving GW 317 then returns a create default bearer response message 329 to the MME 311 that includes the IP address of the WTRU 305. This message 329 also serves as an indication to the MME 311 that the binding was successful. The MME 311 sends an attach accept message 331 to the WTRU 305 through 3GPP access 303. The 3GPP access system 303 initiates radio bearer setup procedures and the 3GPP access system responds with an Attach Complete Message 331. The WTRU 305 may send a binding update 321 to the PDN GW 309 to de-register its DSMIPv6 binding 325 that was created with the WTRU 305 was in untrusted non-3GPP IP access 301. The WTRU 305 may send IKEv2 messages if necessary to tear down its system aspects (SA) with the ePDG 307.\nA handover process 400 from 3GPP IP Access to untrusted non-3GPP IP access over the S2b interface is shown in FIG. 4. The WTRU 401 is connected to the 3GPP network through the serving GW 403 through a PMIPv6 Tunnel 405 to the PDN GW 407 when handover is initiated and the WTRU attaches to the non-3GPP network 411. Authentication of the WTRU 401 is done by the HSS/AAA 409 on the 3GPP system. IKEv2 authorization and tunnel setup 415 between the WTRU 401 and the ePDG 413. A proxy binding update message 417 is sent the PDN GW 407. The proxy binding Acknowledgement message 419 is returned to the ePDG 413. The IPsec tunnel setup and address configuration is performed 421 and the PMIPv6 tunnel 423 is established between the ePDG 413 and the PDN GW 407. Non-3GPP IP access is now established through an IPsec Tunnel 425 between the WTRU 401 and the ePDG 413 and a PMIPv6 tunnel 427 between the ePDG 413 and the PDN GW 407.\nThe handover process 500 from 3GPP IP Access to trusted/untrusted non-3GPP IP access over the S2c interface is shown in FIG. 5. The WTRU 501 is connected to the 3GPP network through the serving GW 503 through a PMIPv6 Tunnel 505 to the PDN GW 507 when handover is initiated and the WTRU attaches to the non-3GPP network 511. Authentication of the WTRU 501 is done by the HSS/AAA 509 on the 3GPP system. IKEv2 authorization and tunnel setup 515 between the WTRU 501 and the ePDG 513 when the non-3GPP network is untrusted. Alternatively, the IKEv2 authorization and tunnel setup 516 may occur between the WTRU 501 and the PDN GW 507 when the non-3GPP network is a trusted network. A proxy binding update message 517 is sent the PDN GW 507. The proxy binding Acknowledgement message 517 is returned to the ePDG 513. The IPsec tunnel setup and address configuration is performed 521 and the PMIPv6 tunnel 527 is established between the WTRU 501 and the PDN GW 507. Non-3GPP IP access is now established through an IPsec Tunnel 525 between the WTRU 501 and the ePDG 513 and a PMIPv6 tunnel 527 between the WTRU 501 and the PDN GW 507."}
{"text": "The processes associated with integrated circuit design generally involve several design transformations from high levels of abstraction to intermediate and lower levels of abstraction. At each level, equivalence verification is performed to prove that any transformation has not functionally or logically altered the design.\nConventional equivalence verification methods generally require a mapping of state elements between the implementations to be verified as equivalent. These mapping methods fall into one of several categories. A naming method, which is fairly common, maps corresponding points in each design that are named similarly. However, the names of state elements may be altered or lost during a design transformation performed by some design tools using the naming method. In addition, the naming method can be fairly unusable for a reverse engineering type application. Another approach is a functional method that makes use of a canonical representation of a state element's fan-in logic cone (e.g., binary decision diagrams or BDDs). However, the functional method can often require that all of the inputs of a cone be mapped when that information might not be available, thereby creating a circular problem.\nAnother approach is a simulation method that simulates both design levels and then matches or narrows the selection of simulation components to state elements with similar value vectors. The simulation method can however be difficult to use in the presence of intentional or unintentional modifications. In another approach referred to as a structural method, a net-list to net-list type comparison is performed that examines combinational cone structure. In some instances, a combination of these approaches is used. However, many or all of the conventional approaches to equivalence verification have some shortfalls. Accordingly, an improved method for performing equivalence verification is needed."}
{"text": "The invention relates to processing multiview video.\nMultiview Video Coding (MVC) relates to compression of video sequences (e.g., a sequence of images or “pictures”) that are typically acquired by respective cameras. The video sequences or “views” can be encoded according to a standard such as MPEG. A picture in a video sequence can represent a full video frame or a field of a video frame. A slice is an independently coded portion of a picture that includes some or all of the macroblocks in the picture, and a macroblock includes blocks of picture elements (or “pixels”).\nThe video sequences can be encoded as a multiview video sequence according to the H.264/AVC codec technology, and many developers are conducting research into amendment of standards to accommodate multiview video sequences.\nThree profiles for supporting specific functions are prescribed in the current H.264 standard. The term “profile” indicates the standardization of technical components for use in the video encoding/decoding algorithms. In other words, the profile is the set of technical components prescribed for decoding a bitstream of a compressed sequence, and may be considered to be a sub-standard. The above-mentioned three profiles are a baseline profile, a main profile, and an extended profile. A variety of functions for the encoder and the decoder have been defined in the H.264 standard, such that the encoder and the decoder can be compatible with the baseline profile, the main profile, and the extended profile respectively.\nThe bitstream for the H.264/AVC standard is structured according to a Video Coding Layer (VCL) for processing the moving-image coding (i.e., the sequence coding), and a Network Abstraction Layer (NAL) associated with a subsystem capable of transmitting/storing encoded information. The output data of the encoding process is VCL data, and is mapped into NAL units before it is transmitted or stored. Each NAL unit includes a Raw Byte Sequence Payload (RBSP) corresponding to either compressed video data or header information.\nThe NAL unit includes a NAL header and a RBSP. The NAL header includes flag information (e.g., nal_ref_idc) and identification (ID) information (e.g., nal_unit_type). The flag information “nal_ref_idc” indicates the presence or absence of a slice used as a reference picture of the NAL unit. The ID information “nal_unit_type” indicates the type of the NAL unit. The RBSP stores compressed original data. An RBSP trailing bit can be added to the last part of the RBSP, such that the length of the RBSP can be represented by a multiple of 8 bits.\nThere are a variety of the NAL units, for example, an Instantaneous Decoding Refresh (IDR) picture, a Sequence Parameter Set (SPS), a Picture Parameter Set (PPS), and Supplemental Enhancement Information (SEI), etc.\nThe standard has generally defined a target product using various profiles and levels, such that the target product can be implemented with appropriate costs. The decoder satisfies a predetermined constraint at a corresponding profile and level.\nThe profile and the level are able to indicate a function or parameter of the decoder, such that they indicate which compressed images can be handled by the decoder. Specific information indicating which one of multiple profiles corresponds to the bitstream can be identified by profile ID information. The profile ID information “profile_idc” provides a flag for identifying a profile associated with the bitstream. The H.264/AVC standard includes three profile identifiers (IDs). If the profile ID information “profile_idc” is set to “66”, the bitstream is based on the baseline profile. If the profile ID information “profile_idc” is set to “77”, the bitstream is based on the main profile. If the profile ID information “profile_idc” is set to “88”, the bitstream is based on the extended profile. The above-mentioned “profile_idc” information may be contained in the SPS (Sequence Parameter Set), for example."}
{"text": "Semiconductor memory devices are generally organized in a bidimensional array, or memory matrix, wherein the memory elements are located at the intersection of rows, or \"word lines,\" and columns, or \"bit lines,\" of the matrix. To access a given memory element, it is necessary to select the word line and the bit line at the intersection of which the memory element is located. To this purpose, the memory address bus is divided into row and column address signals, which are decoded independently.\nIn the manufacture of semiconductor memories, defects are frequently encountered that afflict a limited number of memory elements in the memory matrix. The reason for the high probability of defects of this type resides in that in a semiconductor memory device the greatest part of the chip area is occupied by the memory matrix. Moreover, it is in the memory matrix, and not in the peripheral circuitry, that the manufacturing process characteristics are usually pushed to limits.\nIn order to prevent rejection of an entire chip due to the presence of a limited number of defective memory elements, out of many millions of memory elements, and therefore increase the manufacturing process yield, the typical manufacturing technique provides for a certain number of additional memory elements, commonly called redundancy memory elements. Redundancy memory elements are to be used as a replacement of those elements that, during testing of the memory device, prove defective. The selection circuits, with which the integrated component must necessarily be provided, and which allow the above-mentioned functional replacement of a defective memory element with a redundancy memory element, are indicated as a whole with the name of \"redundancy circuitry,\" while the set of redundancy memory elements and circuitry is defined as \"redundancy.\"\nThe redundancy circuitry comprises programmable non-volatile memory registers, or redundancy registers, suitable to store those address configurations corresponding to the defective memory elements. Such registers are programmed once during the memory device testing, and must retain the information stored therein even in absence of the power supply.\nIn practical implementations of redundancy in memory devices, both word lines and bit lines of redundancy memory elements are generally provided in the memory matrix. Each redundancy word line or bit line is associated with a respective row or column redundancy register, wherein the address of a defective word line or bit line is stored. Whenever the defective word line or bit line is addressed, the corresponding redundancy word line or bit line is selected.\nAs far as word lines are concerned, it has been recognized that the most frequent defect consists in short-circuits between adjacent word lines. This situation is, however, easily detected during testing. When the selection of one of two short-circuited word lines is attempted, the potential of such word line cannot rise to the designed value, being linked by the short-circuit to the potential of the adjacent non-selected word line. Therefore, when a defective word line is found during testing, it is assumed that such word line is short-circuited with the adjacent word line and both the word lines are replaced by two respective redundancy word lines. From then on, the two defective word lines will never be selected. The adjacent word line, in the scanning sequence, follows the defective word line.\nSince defective word lines always come in pairs, it is known to design the row redundancy registers in such a way as to store a pair of row addresses in each row redundancy register. Therefore, each row redundancy register is associated with a respective pair of redundancy word lines: into each row redundancy register can therefore be programmed the addresses of two adjacent short-circuited word lines.\nEach row redundancy register comprises programmable non-volatile memory cells wherein the addresses of two adjacent defective word lines can be programmed. Each one of such memory cells must comprise at least one programmable non-volatile memory element, such as a fuse or a floating-gate MOSFET, a load circuit for reading the information stored therein, and a program load circuit for the programming of the memory element according to the logic state of a respective address bit in the row address signal set, or row address bus.\nIn memory devices the association between the word lines in the memory matrix and the respective selection signals generated by the row address decoding and selection circuitry is generally such that adjacent word lines have addresses which only differ in one bit. Since, however, such bit can be any one of the bits constituting the row address signal set, it follows that in order to be always able to replace two adjacent short-circuited word lines whatsoever, each row redundancy register should store two full row addresses. This means that each row redundancy register must comprise a number of memory cells equal to twice the number of row address bits. Since the memory cells occupy each a significant chip area, this would lead to an excessive increase in the overall chip size, so that the overall process yield is decreased instead of increased. This is why a compromise is generally reached between the repairability rate for defective word line pairs and the increase in chip size. In practice, designers give up the possibility to replace two adjacent short-circuited word lines whatsoever, limiting such replacement to adjacent word lines whose addresses only differ in one or more bit belonging to a given subset of the whole row address signal set. If, for example, the row address signal set comprises m address bits, it can be thought of as being the sum of two subsets n and q such that n contains the most significant row address bits, while q contains the least significant row address bits. Limiting the replacement of defective word line pairs to take place for adjacent word lines having addresses differing in one or more bits of the subset q, it is sufficient to store in a given row redundancy register the full row address m for one of the pair of adjacent word lines, and the subset q for the other word line of the pair. This means that each row redundancy register must be made up of m+q memory cells, instead of 2 m. The impact on the repairability rate can be appreciated by considering that the probability of having a short-circuit defect between two adjacent word lines with addresses differing in one or more bits in the subset n is 1/2 q.\nIn the European Patent Application No. 93830474.8, incorporated herein by reference provided as background and not admitted as prior art, a program load circuit for the programming of a memory cell in a non-volatile memory register, such as a row redundancy register, is described. The datum to be programmed into the memory element of the memory cell can be directly supplied by one of the address signal lines already present in the memory device for supplying the decoding circuitry, without the need of generating additional signals."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an auto-contrast enhancement system, and more particularly to a human visual system (HVS)-based auto-contrast enhancement system with side-effect suppression.\n2. Description of Related Art\nImage enhancement is widely applied to refine images with low contrast or images captured in non-ideal lighting conditions. Traditional contrast enhancement algorithms usually suffer from over-enhancement, which makes images look unnatural and amplifies noise.\nHistogram enhancement (HE) is the most common method for improving contrast enhancement. With HE, the intensities of an image are distributed more uniformly to achieve better contrast. The overall complexity of HE-based algorithms is low compared to that of other methods. However, HE can easily induce over-enhancement, noise amplification, fine detail loss, or making the image look unnatural.\nMoreover, human eyes cannot easily distinguish the intensity variation in a saturated region. The saturation region is generally regarded as a smooth region that contains no texture information. However, slight intensity variation and random distributed hue are frequently hidden inside the saturation region. If contrast enhancement is applied directly, noise amplification will be strong in this region.\nTherefore, a proper enhancement level and side-effect reduction methods are crucial in a contrast enhancement system."}
{"text": "This invention relates to a damper for a motor vehicle suspension system that includes a pneumatic self-leveling capability to maintain the ride height of the vehicle at a desired level.\nMotor vehicles, such as passenger cars and light trucks, experience wide variations in loading. Accordingly, it is desirable to adjust the spring rate and load of the suspension system to compensate for variations in vehicle loading, and accordingly maintain the ride height of the vehicle despite the variations in loading.\nSeveral devices are known in the prior art that compensate for varying vehicle loads to maintain ride height. These devices include air springs mounted on the vehicle either parallel to the main spring or as a replacement for the main spring, and an air spring in combination with a suspension damper, commonly referred to as an xe2x80x9cair shockxe2x80x9d. All of these devices require an air compressor to provide compressed air to the device, a sensor which measures ride height and actuates the compressor, and a solenoid valve to release air when the vehicle loading is reduced. Accordingly, these devices are expensive, difficult to install, and require that the vehicle have the necessary electrical and mechanical infrastructure to permit their use. Hydro-pneumatic systems, in which hydraulic fluid is pumped into the gas chamber of a gas charged damper to adjust ride height are also used to maintain ride height. These systems also require pumps, valves, etc. and tend to be costly and heavy. Finally, a self-contained device sold by Mannesmann Sachs AG of Schweinfurt, Germany under the trade name Nivomat includes a pumping mechanism, level sensor, hydro-pneumatic spring and damper that uses normal road undulations to pump up the suspension to maintain ride height. This device is easy to install but is expensive and compromises the ride quality of the vehicle suspension.\nAccording to the present invention, a suspension damper includes a reservoir housing with a piston slidably mounted therein and including a piston rod extending from the reservoir housing and attached to the vehicle by a mount assembly. An air chamber is defined by a sleeve or dust cover cirucmscribing the piston rod and connected to the reservoir housing by a flexible sleeve. The mount assembly includes a pumping chamber which expands and contracts in response to normal road undulations to inflate the air chamber. A valve controlled by the flexible sleeve controls communication from the air chamber to ambient atmosphere to thereby control inflation of the air chamber, thereby controlling ride height."}
{"text": "1. Field of the Invention\nThe present invention relates to a luminescent storage screen of the type suitable for latently storing an x-ray image in an x-ray diagnostics installation, and an apparatus for reading-out the stored image from the screen.\n2. Description of the Prior Art\nAn x-ray diagnostics installation is described in German OS No. 35 29 306 which includes a luminescent storage screen having a photo stimulable phosphor which is irradiated by x-rays attenuated by an examination subject. The phosphor serves as a radiation-sensitive transducer in which holes (in the electronic sense) are generated, and are stored in potential traps or wells. The entire surface of this storage screen is scanned pixel-by-pixel during read-out by a different radiation source, for example a laser, so that the holes stored in the traps are excited, and then fall back into a stable energy level. The energy difference is emitted in the form of light quanta. The x-ray image latently stored in the screen can thus be read-out of the screen at a later time.\nThe storage screen in this known installation consists of two phosphors separated by a light-impermeable layer. These phosphor layers are scanned surface-wide by a beam deflected by a mirror, for example, in the vertical direction, with the beam being divided by a beam splitter so that both phosphor layers are simultaneously scanned. All picture elements in a line of the storage screen are successively scanned by this deflection. The entire screen surface is scanned by moving the storage screen perpendicularly relative to the line scanning, advancing the screen line-by-line. The light emitted by the stimulable phosphor layers as a result of the interaction with the read-out beam is acquired by two optical collectors, and is supplied to the light-sensitive input faces of two separate light detectors. The respective output signals of the detectors are supplied to a video chain for reproduction on a display, with the output signals of the two detectors being combined.\nInstead of dividing the read-out beam, however, it is also possible to use two separate radiation generators, each undertaking a separate scanning of one of the phosphor layers. In this case, however, it is necessary to synchronize the deflection of both scan beams, so that the picture elements of the two phosphor layers which correspond to each other are simultaneously read-out. Synchronization need not be undertaken at the time of read-out; it can alternatively be undertaken electronically in a processing circuit following the light detectors.\nThe two storage layers in combination function as a single storage layer having an increased thickness. It is known that the efficiency of a phosphor layer increases with increasing thickness, since more x-ray radiation is absorbed in a thicker layer. Image sharpness, however, decreases with increasing layer thickness, so that fine details of the image can no longer be adequately portrayed. By using two separate phosphor layers, with a separate read-out of each layer, the efficiency of a layer having a thickness corresponding substantially to the combined thickness of the two layers is achieved, while a better image sharpness is obtained than would be possible using a single layer of the same combined thickness. A disadvantage is the complexity of the read-out system which is required for such a screen, which not only requires two separate read-out beams, but also structure for guiding and synchronizing the two beams so that the pixels separately read-out from each phosphor layer can be correlated to form a single image. Two synchronized scan systems, or at least two synchronized scan beams, are required.\nAnother x-ray diagnostics installation is described in OS No. 29 51 501 which includes a luminescent storage screen having a phosphor layer scanned by a read-out beam, with detectors being disposed on both sides of the storage screen. These detectors respectively collect the light emitted from the phosphor layer in an upward direction and the light emitted in a downward direction. A higher light yield and thus an improved bright-dark contrast are thereby achieved. Because only a single phosphor layer is used, however, the competing problems of efficiency of x-ray radiation absorption and image sharpness are still present."}
{"text": "Notebook (also called laptop) computers are lightweight personal computers, which are quickly gaining popularity. The popularity of the notebook computers has especially increased since their prices have been dropping significantly, while maintaining similar performance as their larger siblings (i.e., desktop computers or workstations).\nOne clear advantage of notebook computers is their ease of portability. The lighter weight and increased battery life has made the use of notebook computers more commonplace. Aside from size, a major difference between notebook computers and personal computers is the graphical display system. Notebook computer designers employ a variety of techniques to produce lightweight graphical display systems.\nFIG. 1 shows a typical notebook computer in accordance with the prior art. A notebook computer 100 typically includes a liquid crystal display (LCD) 102 for displaying graphical images and a base unit 104 for encasing different hardware elements of the notebook computer. The hardware elements include a keyboard 106, a battery 108, a hard drive (not shown), and a central processing unit (CPU, not shown).\nIn some applications, the notebook computers may be utilized as a tablet (also called a convertible laptop), i.e., where a user may input data by inputting it directly onto the LCD 102. Some users, however, dislike tablet computers because their usage requires a different approach (e.g., using a pen instead of a pointing stick). Additionally, the current convertible laptops use their own proprietary or nonstandard operating systems (OSes). The nonstandard implementation of the convertible laptops creates compatibility issues with other computer systems in use at the user's home or office.\nThe wider use of notebook computers has also created problems for their users, including physical ones. Incidents of carpal tunnel syndrome, eyestrain, and backache are on the rise. For example, increased passenger capacity in airplanes has resulted in reduced leg and shoulder room, in turn, increasing the stress on laptop users' posture when utilizing their computers. For example, with respect to FIG. 1, users often rest their palms on the hand rest area of the keyboard (210) while typing on the keyboard 106, which increases the stress on the user's hands. Similarly, user eyestrain may result when the LCD panel 102 cannot be tilted back far enough (e.g., due to space restrictions), thereby rendering a less than ideal brightness yield from the incorrectly angled LCD panel."}
{"text": "The number of new head and neck cancer incidences around the world is approximately 600,000+/− per year, with approximately 40,000+/− being within the US. Decreased or absent output from the salivary glands and the cells of membranous linings found within the mouth, throat and airway passages has been strongly linked, but is not limited to association with cancer treatments such as radiation and/or some chemotherapy agents, especially those treatments involving head and neck cancers. There is also a significant population of survivors who suffer cancer treatment/radiation induced xerostomia of a permanent nature. Xerostomia (dry mouth) is one of the most widely cited manifestations associated with cancer treatments of the head and neck resulting in decreased salivary gland output.\nSaliva production typically begins to decrease very soon after the start of head and/or neck radiation therapy, often within the very first week. This reduction in salivary function continues declining as the treatment progresses. The extent of damage that occurs to the moisture producing capabilities of the salivary glands and functioning of the cells of the membranous linings that lubricate the throat and airways will vary, dependent upon the exact location of the area treated as well as the total dose of radiation that was delivered.\nDry mouth can also be caused by activity, exertion and certain medications or diseases such as Sjrojen's. While this type of dry mouth can be temporary, radiation-therapy-induced xerostomia can in some cases completely destroy salivary gland function, has no known cure at this time and is considered permanent after a certain point in recovery. This means that patients who received treatment for head, neck or chest related cancers likely have an affected ability to produce moisture within the oral cavity and will likely suffer with this condition for the remainder of their lives.\nThere are a wide variety of bio-engineered, herbal-based and synthetic products and solutions currently available for relief of symptoms of xerostomia. Some of these products may operate through simulation of salivary excretions and some through stimulation of remaining working salivary glands. These products include gums, mouthwashes, sprays, and toothpastes, etc. . . . . While these solutions may offer some relief there are still problems. For example many of these are effective for only short periods, may have limitations upon number of times recommended to be used per day, must be carried and/or applied, may have taste, textures, an oily or otherwise unnatural feel in the mouth. Additionally, the increase or introduction of actual environmental moisture delivered is generally minimal. And, the delivery of the solutions can fail to reach all aspects of oral cavity and may not provide sufficient hydrating lubrication in the form of moisture that is delivered to the membranous linings of airways. Additionally, severe xerostomia can result in water overload causing feelings of stomach fullness and/or may even be the cause of hyponatremia. It would therefore be desirable to have an oral hydrating system that is convenient, and may be used as often as user wishes, as many times per day without limit as needed."}
{"text": "This invention relates generally to acoustic signal correlation, and more particularly to optical correlation using a light emitting diode (LED) array.\nBroadband acoustic processing attempts to detect a target and track its path by the use of the target's wide frequency band energy sources. One important processing technique suited to broadband signal detection is correlation. Cross correlation is done by processing the signals from two spatially separated transducers or sensors, such as from two sonobuoys or from the multihydrophone array of a single buoy. The channels of information from these two transducers contain common broadband signals, which, because of the transducer spacing, are detected apart in time. The time delay depends upon the location of the target with respect to the two transducers. The correlator detects the broadband signals and determines the time delay. Correlation done in this way over time allows the target to be tracked.\nBroadband acoustic processing requires the ability to perform intra-/inter-sensor data correlation at high speeds. Currently, this operation is performed digitally with high-speed array processing systems using both time domain and frequency domain methods. These processors make use of the fast Fourier transform and the inverse fast Fourier transform algorithms. Although these digital processors are fast, they are also large, complex, and consume large amounts of power.\nThe invention disclosed herein is an optical correlator. The signals from the transducers are processed optically, using a light emitting diode (LED) array and a charged coupled device (CCD) imaging array. Currently available optical correlators use acousto-optical cells having a fixed delay range equal to the transit time through the cell. These correlators are too fast to be used in low data rate processing and lack signal delay range flexibility and sensitivity."}
{"text": "This invention relates to method and apparatus for decoding PCM digital input words, and more particularly to an integrated PCM decoder which is substantially insensitive to parasitic and stray capacitance effects.\nPCM encoders and/or decoders are described in the articles \"A Two-Chip PCM Voice Codec With Filters\" by Y. A. Haque, et al., IEEE Journal of Solid State Circuits, Vol. SC-14, No. 6, pages 961-969, December 1979; \"A PCM Voice Codec With On-Chip Filters\" by J. T. Caves, et al., IEEE Journal of Solid State Circuits, Vol. SC-14, No. 1, pages 65-73, February 1979; \"A Segmented Mu-255 Law PCM Encoder Utilizing NMOS Technology\" by Paul R. Gray, et al., IEEE Journal of Solid State Circuits, Vol. SC-11, No. 6, pages 740-747, December 1976; and \"A Unified Formulation of Segment Companding Laws and Synthesis of Codecs and Digital Compandors\" by H. Kaneko, The Bell System Technical Journal (BSTJ), September 1970, pages 1555-1588, which are incorporated herein by reference. Mu-law decoders operating on 8-bit PCM code words have employed arrays of binary weighted capacitors, the total array capacitance typically being 255 times the capacitance of the smallest capacitance of the array. The capacitance of the two largest binary weighted capacitors of the array are therefore 128 and 64 times that of the smallest capacitor thereof. In integrated circuit structures it is desirable to minimize the surface area of a chip that is required for individual circuit elements. It is readily seen that the chip area dedicated for integrated capacitors in such an integrated decoder may be reduced by approximately 50% by deleting only one of the binary weighted capacitors. An integrated decoder requiring a capacitance of only 32 times that of the smallest capacitor of an array is disclosed in copending patent application (infra).\nAn object of this invention is the provision of an improved PCM decoder requiring only three equal valued capacitors."}
{"text": "The present invention relates generally to a system and method for providing a user with a plurality of guaranteed minimum retirement income payments and, more particularly, to a system and method for guaranteeing minimum periodic retirement income payments using a defined benefit account inside of a defined contribution plan.\nUp until about 1870, more than half of the United States' adult workers were farmers. These adult workers were typically engaged in their occupations until their death or until their health prevented them from continuing their occupations. It was uncommon to have a prolonged retirement period before a worker's death.\nAfter 1870, however, industry developed rapidly and the economy tended increasingly to be characterized by industrialization and urbanization. The result was that workers increasingly were employed in more industry-related jobs and became more dependent upon a continuing flow of monetary income to provide for themselves and their families. Additionally, the average life expectancies of workers began to increase significantly. It became more common for workers to retire from employment and to survive for longer periods of time following their retirements. Retirement programs began to take hold. The Social Security program was introduced in 1935 and had an old age insurance component which provided a lump sum benefit for workers at age 65. At that time, the average life expectancy of a worker was 68.\nCurrently, however, half of male workers reaching age 65 can expect to still be alive at age 82 and half of female workers reaching age 65 can expect to be alive at age 86. The Social Security program is not keeping pace with such changes. The number of employees entering the workforce has been less than the number of new retirees for the last several years and this trend is expected to increase as the “Baby Boomers” age. The Social Security Administration (“SSA”) projects a shortfall in its trust fund which provides benefits to retirees beginning in 2013. The SSA believes that an immediate and permanent increase of social security payroll taxes is necessary in order to enable it to pay for the full amount of old age benefits it currently provides retirees. Now, employees and employers contribute approximately 12.4 percent of salaries to the Social Security trust fund. The SSA projects that contributions must be increased to at least 38 percent in order for its trust fund to remain fully funded. Therefore, it is becoming increasingly uncertain whether the Social Security program will continue to remain viable until the time that today's workers are ready to retire. Moreover, many retirees have found that the amount of retirement benefits to which they are entitled under the Social Security program is insufficient to enable them to maintain a desired level of comfort in their retirement. They have found a need to supplement such Social Security benefits with income from other sources.\nIn addition to the institution of the Social Security program in the 1930s, beginning in the early 1900s, it became increasingly more common for employers to provide their workers, or employees, with some sort of retirement benefits or pensions. These retirement benefits or pensions were originally designed, in part, to reward an employee for his/her long career with a company and to help provide an income once such employee retired. Such retirement benefits or pension plans therefore required minimum periods of employment before an employee's entitlement to the pension amount became vested. However, many such retirement benefits or pensions are not portable. In other words, if an employee leaves the employ of an employer, that employee may lose all entitlement to such retirement benefit or pension if the employee terminates his/her employment prior to the expiration of the vesting period. This was not a problem when employers first instituted such retirement benefits or pension plans as employees tended to remain employed with one employer for their entire career until they retired.\nHowever, in today's mobile society, employees do not tend to remain employed by one employer for their entire careers. Many employees therefore lose some or all of their projected retirement benefits which may have accrued during their employ by their employers when they leave the employ of such employers.\nFurthermore, in addition to the trend of a more mobile society and an increased level of employment changes, many employers are decreasing the numbers of their employees and are instead increasingly turning to non-employee labor in part to cut expenses resulting from employee benefits such as costs related to funding employee retirement plans. Thus, many individuals in the workforce today are technically not considered “employees” but instead are independent contractors for whom employment benefits such as retirement benefits are not provided. Additionally, many employers are ceasing to offer defined benefit plans altogether because of the costs. In fact, according to statistics published by the Pension Benefit Guaranty Corporation, defined benefit pension plans of employers have decreased by more than 60 percent since 1985, with the number of U.S.-based employers that offer such defined benefit pension plans decreasing from 114,000 in 1985 to less than 40,000 in 1999. Only 21.3 percent of working family heads are currently covered by an employer-funded defined retirement benefit or pension plan.\nBecause of the decrease in the number of employers that offer defined retirement benefit pension plans, the decrease in the number of workers entitled to employer-funded retirement benefits and also because of the increased mobility of the workforce resulting in the loss of such employer-funded benefits, many workers have started to fund their own retirement savings plans. Tax laws have enabled workers to realize tax benefits from deferring their income by putting amounts from their paychecks into such retirement savings plans. Increasingly, such employee-self-funded retirement savings plans are becoming the primary sources of income on which employees survive following retirement.\nHowever, one disadvantage of the increased reliance upon employee-self-funded retirement savings is that these plans do not provide a level of retirement income that is guaranteed for the employee. In addition, many employees do not have any idea of an amount required to be saved in order to achieve a desired level of income to ensure a comfortable lifestyle upon their retirement. Thus, they do not contribute a sufficient amount of their salaries towards such retirement savings to provide an adequate income level to maintain the standard of living they desire upon retirement. Based on the results of the Retirement Confidence Survey sponsored by the Employee Benefits Research Institute (EBRI), the American Savings Education Council (ASEC), and Matthew Greenwald and Associates, 22 percent of all employed adult workers have saved less than $10,000 towards retirement, 50 percent have saved less than $50,000 and only 25 percent of adult workers over the age of 55 have accumulated more than $100,000.\nRetirement income needs may increase in the event such retirees suffer from health-related problems. In fact, many employees today express concern that they will not have adequate funds saved to provide for themselves during their retirement in the event they suffer health-related problems after they retire. They are currently seeking some means to ensure a higher level of income saved for such crises.\nEmployees often do not participate in their employer-sponsored retirement savings plans which will increase the level of their savings through interest income or a return on investment. Also, many individuals lack the sophistication needed to determine the appropriate type of investment vehicle which will offer them a high return on their investment but which is also secure enough so that their savings are not placed at risk by a high-risk type of investment vehicle.\nThus, there is a need for an investment vehicle which will provide a future minimum retirement income and a current account value, both of which are portable so that a worker will not lose any value or income vested in a fully funded investment vehicle if the worker leaves the employ of an employer or changes jobs.\nThere is also a need to provide a defined retirement benefit which will guarantee an individual a minimum defined income level upon the individual's retirement.\nAdditionally, there is a need for a retirement investment vehicle which may provide a guaranteed minimum level of retirement income and also may afford an individual an opportunity for an increase in value of the benefits provided if market performance of the retirement vehicle exceeds a predefined benchmark.\nAdditionally, there is a need for an investment vehicle that addresses multiples of these concerns and may help workers accumulate additional wealth, and more tangibly plan for income generation during retirement."}
{"text": "Limited oil resources and rapidly increasing crude oil prices have created a need for lubricants with longer useful lives. Also, longer intervals between crankcase oil changes will reduce the volume of used oil for disposal. For these and other reasons, the efficiency and useful lives of oil-based lubricants, particularly crankcase lubricants, must be improved.\nOxidation of the oil component in the lubricant substantially shortens its useful life. Oxidation yields corrosive acids and an undesirable increase in viscosity. While high quality basestocks tend to be relatively resistant to oxidation, contaminants (e.g., iron) and common additives can greatly accelerate oxidation. Inclusion of detergents (e.g. calcium or magnesium detergents) and dispersants (e.g. polyamine or polyester derivatives of alkenyl succinic acids or anhydrides) is desirable for oil performance, but these additives accelerate oxidation to such an extent that oxidation is a major cause of reduced useful life.\nWith depletion of high quality basestock oil reserves, resort to lower quality basestocks has become necessary. These lower quality basestocks have a greater tendency to oxidize than do higher quality basestocks.\nIf the life of a crankcase lubricant is to be maximized, oxidation must be minimized. Over the years various oxidation inhibitors, or antioxidants, have been proposed. Examples of antioxidants which have been proposed for use in crankcase lubricants include zinc dihydrocarbyl dithiophosphates which are primarily used as antiwear agents but also act as antioxidants, aromatic amines (e.g. alkylated phenylamines and phenyl-.alpha.-naphthylamines), hindered phenols, alkaline earth metal salts or sulfurized alkyl phenols in which the alkyl groups have 5 to 12 carbon atoms (e.g., calcium nonylphenyl sulfide and barium octylphenyl sulfide), phosphosulfurized or sulfurized hydrocarbons, and oil soluble copper compounds.\nSome of the above mentioned antioxidants are very effective. Thus, European Patent No. 24 146 B teaches lubricating compositions comprising a major amount of lubricating oil, from 1 to 10 wt % of certain ashless dispersants or from 0.3 to 10 wt % of certain nitrogen- or ester- containing polymeric viscosity index improver dispersants, or mixtures of dispersants and viscosity index improver dispersants; 0.01 to 5 wt % of zinc dihydrocarbyl dithiophosphate (ZDDP); and 5 to 500 parts per million (ppm) by weight of added copper in the form of oil-soluble copper compound. The patent notes that the inexpensive copper antioxidants are effective at low concentrations and therefore do not add much to the products cost. In the amounts employed, the copper compounds do not interfere with the performance of other components of the lubricating system. In many instances, completely satisfactory results are obtained when the copper compound is the sole antioxidant in addition to ZDDP. Alternatively, the patent teaches that for particularly severe conditions where a supplementary antioxidant may be desirable, the amount of supplementary antioxidant required is small, often far less than the amount required in the absence of the copper compound. Supplementary antioxidants mentioned include phenols, hindered phenols, bis-phenols, sulfurized phenols, catechol, alkylated catechols, sulfurized alkyl catechols, diphenylamine, alkylated diphenylamines, phenyl-l-naphthylamine and its alkylated derivatives, alkyl borates, aryl borates, alkyl phosphates, aryl phosphites, aryl phosphates, O,O,S-trialkyl dithiophosphates, O,O,S-triaryl dithiophosphates, and O,O,S-trisubstituted dithiophosphates containing both alkyl and aryl groups.\nA specific copper-containing lubricant described in U.S. Pat. No. 4,705,641 comprises (A) a basestock and (3) a copper salt and a molybdenum salt wherein the total concentration of the copper and molybdenum metal or metal ions in solution ranges between about 0.006 and about 0.5 wt % of the basestock (60 to 5000 ppm by weight). The stated preferred concentration of copper and molybdenum ranges from about 0.009 to about 0.1 wt % of the basestock (90 to 1000 ppm by weight).\nEP 280,579 A and EP 280,580 A each mention in comparative Example 4 lubricating compositions containing copper oleate and molybdenum oleate.\nDespite the beneficial effect of using soluble copper described in EP 24,146 B, the antioxidant literature contains many suggestions that copper should be avoided. For example, an antioxidant system comprising in association (a) a particular sulfur-containing molybdenum complex and (b) an aromatic amine is described in UK patent 2,097,422. The recited sulfur containing molybdenum complex is prepared by reacting an acidic molybdenum compound with a basic nitrogen-containing substance and a sulfur source. Preferably the reaction is carried out in the presence of a polar promoter. Example 10 describes a test for oxidation stability wherein copper is used as an oxidation catalyst.\nAnother patent teaching that copper is an oxidation catalyst is EP 404,650 A. That application discloses oil-soluble overbased additives comprising an alkali or alkaline earth metal carbonate in combination with a substantially hydrocarbon insoluble organic molybdenum derivative. The application states that the molybdenum derivative is solubilized during the overbasing reaction, perhaps by incorporating the derivative into the micelles of the metal carbonate colloid. The molybdenum complex may be a molybdenum-amine complex formed by reacting an acidic molybdenum compound with an amine, e.g. a primary aliphatic amine. Also taught are oxygen containing molybdenum complexes formed by reacting a molybdenum compound with an oxygen containing compound, e.g. a glycol. Example 16 describes a TFOUT (thin film oxygen uptake test) using a naphthenate of lead, copper, iron, manganese and tin as catalyst.\nDespite the breadth of the antioxidant literature, highly effective antioxidants for lubricating compositions, particularly crankcase lubricants, in which the proportion of metal containing antioxidants can, if desired, be kept low, are still needed."}
{"text": "Processes for stuffing meat products have been continually improved over the years.\nU.S. Pat. No. 4,155,212, for example, relates to a mechanism for wrapping a netting around a meat piece, provided with a cylinder around which is placed the netting in which the meat piece supplied from a table is introduced manually and pushed manually or by a piston. The meat is pushed through the inside of the cylinder and leaves out of its discharge opening wrapped in the netting.\nU.S. Pat. No. 4,621,482 considers the additional application of a collagen film that, together with the netting, wraps the meat in a stuffing apparatus.\nAn evolution of these systems is found in U.S. Pat. Nos. 4,910,034 and 4,958,477, which respectively show a process and an apparatus for production of meat products. The process for obtaining the final product consists mainly of simultaneously wrapping the meat products with a collagen film and an elastic netting supplied from tubes concentric to the central tube in which the meat pieces travel. Both the meat feeding and impulsion through the tube are manual in these systems.\nThe manual stuffing process generally requires one person to feed the meat pieces to the machine and another person to receive the stuffed product from the applicator tube, conform and clip it and deliver it to a third person who clips it. The working speed of this system is not greater than 5 pieces per minute. Applicators exist that are connected to a stuffing machine that feeds the meat, but their speed is the same, although they work with one less person.\nThe TCM machine of Tipper Tie incorporates a pneumatic piston for pushing the meat and an automatic clipping system at the applicator outlet that allows a rate of 7-8 pieces per minute. However, it has limitations regarding the dimensions of the clipping system, its slowness and the need to feed the machine manually interrupting the process during loading.\nPatent application WO 03/066440 describes a method and apparatus for automatic stuffing of meat pieces in a double casing of film and netting, in which participate a synchronised stuffing machine and fast clipper and where the meat piece is made to pass through a tube with a small diameter independent of but smaller that the size of the finished product. The rates reached depend on the flow of meat that the stuffing machine can impel and the time taken by the clipper to clip, reaching a production on the order of 30 pieces per minute.\nThe machine that constitutes the object of the aforementioned patent application WO 03/066440 is conceived to stuff only meat made of small pieces that can be fed by pumping by a stuffing machine.\nThere are, however, a number of products that must be manufactured without pumping. These are those made of entire muscles. These are top-end products that require the integral maintenance of the muscle structure and those in which it is desired to have the product wrapped with its skin visible on the surface. These products cannot be stuffed by a pumping means nor through small diameter tubes, as they would lose their structure and order.\nIf the apparatus described in the aforementioned U.S. Pat. No. 4,958,477 is provided with a mechanical impulsion means, such as a pneumatic piston, and combined with an “iris” type clamp clipper such as that described in U.S. Pat. No. 6,705,063, which allows passage of pieces with a diameter of up to 180 mm and a high-rate clipping, this may provide a first step towards mechanisation of the stuffing process of said products that cannot be pumped. However, to improve the results obtained by the aforementioned TCM machine of Tipper Tie, this system would require an automatic stuffing system that reduces to a minimum the loading time of the entire muscles.\nIn manual loading systems the meat pieces are loaded in a chamber of the applicator after the cylinder has retracted from its pushing stroke. To do so, the chamber lid must be opened and the meat piece must be placed in it orderly and with a specific orientation, taking care that the skin, if it must be visible from the outside in the stuffed product, is properly extended in the surface. The preparation time involved in positioning the meat piece and the subsequent closing of the lid, whether manually or mechanically, until the meat is finally impelled by the piston, implies lost time for the system that will remain unproductive.\nThe automation of the feeding and its operative optimisation therefore constitute the objectives of the invention described below."}
{"text": "1. Field\nThe present invention relates generally to electrical power supplies and, more specifically to uninterruptible power supplies.\n2. Description of the Related Art\nComputer racks, such as server racks, are generally used to house and in some cases interconnect collections of computing devices, like servers and associated storage, power supplies, network switches, and the like. In many cases, the computing devices are relatively numerous and arranged in a relatively high-density array due to the cost of space appropriate to store such computing devices and the desire to reduce latency by having the devices close to one another.\nOften, the computing devices run computing applications having higher up-time goals than can be provided by grid electrical power. Power outages can cause data centers to cease operations. In some cases, backup power supplies, such as diesel generators and batteries are used to supply backup power in the event of a power outage. But in many cases, these systems are too slow to respond to an outage (e.g., many diesel generators), or are larger than is desired for modular data center rack designs, where the backup power supply is built into a rack (e.g., many battery-based backup power supplies)."}
{"text": "1. Field of the Invention\nEmbodiments of the invention relate to an image display device for displaying a two-dimensional plane image (hereinafter referred to as ‘2D image’) and a three-dimensional stereoscopic image (hereinafter referred to as ‘3D image’).\n2. Discussion of the Related Art\nAn image display device displays a 3D image using a stereoscopic technique or an autostereoscopic technique.\nThe stereoscopic technique, which uses a parallax image between left and right eyes of a user with a high stereoscopic effect, includes a glasses type method and a non-glasses type method, both of which have been put to practical use. In the glasses type method, the parallax image between the left and right eyes is displayed on a direct-view display or a projector through a change in a polarization direction of the left and right parallax image or in a time-division manner, and thus a stereoscopic image is implemented using polarization glasses or liquid crystal shutter glasses. In the non-glasses type method, an optical plate such as a parallax barrier for separating an optical axis of the parallax image between the left and right eyes is generally installed in front of or behind a display screen.\nAs shown in FIG. 1, the image display device using the glasses type method may include a patterned retarder 5 for converting polarization characteristics of light incident on polarization glasses 6 on a display panel 3. In the glasses type method, a left eye image (L) and a right eye image (R) are alternately displayed on the display panel 3, and the polarization characteristics of light incident on the polarization glasses 6 are converted by the patterned retarder 5. Through this operation, the glasses type method implements a 3D image by spatially dividing the left eye image (L) and the right eye image (R). In FIG. 1, a reference numeral 1 denotes a backlight unit providing light to the display panel 3, and reference numerals 2 and 4 denote polarizing plates respectively attached on upper and lower surfaces of the display panel 3 so as to select a linear polarization.\nIn the glasses type method, visibility of the 3D image is degraded due to crosstalk generated at the position of an upward or downward viewing angle. As a result, in the general glasses type method, the upward/downward viewing angle capable of allowing the user to view the 3D image of the good image quality is very narrow. The crosstalk is generated because the left eye image (L) passes through a right eye patterned retarder region as well as a left eye patterned retarder region and the right eye image (R) passes through the left eye patterned retarder region as well as the right eye patterned retarder region at the position of the upward/downward viewing angle. Thus, as shown in FIG. 2, Japanese Laid Open Publication No. 2002-185983 discloses a method for obtaining a wider upward/downward viewing angle by forming black stripes (BS) in patterned retarder regions corresponding to black matrixes (BM) of a display panel to thereby improve the visibility of the 3D image. In FIG. 2, when observing at a predetermined distance (D), a viewing angle (α), at which the crosstalk is not theoretically generated, depends on the size of black matrixes (BM) of the display panel, the size of black stripes (BS) of the patterned retarder, and a spacer (S) between the display panel and the patterned retarder. The viewing angle (α) widens as the size of the black matrixes and the size of the black stripes increase and as the spacer (S) between the display panel and the patterned retarder decreases.\nHowever, the related art image display device has the following problems.\nFirst, the black stripes of the patterned retarder used to improve the visibility of the 3D image through the improvement of the viewing angle interact with the black matrixes of the display panel, thereby generating moiré. When a 2D image is displayed, the visibility of the 2D image is much degraded. FIG. 3 shows the results obtained by observing a 47-inch display device sample at a location 4 meters away from the display device to which the black stripes are applied. When the 2D image is displayed, moirés of 90 mm, 150 mm, and 355 mm are visible based on observation positions A, B, and C, respectively.\nSecond, the black stripes used to improve the visibility of the 3D image through the improvement of the viewing angle bring about a side effect allowing a luminance of the 2D image to be drastically degraded. This is because, as shown in FIG. 4(b), in the related art, predetermined portions of pixels of the display panel are covered by the black stripe patterns. Accordingly, when the 2D image is displayed, an amount of transmitted light is reduced by about 30% compared with the case where the black strips are not formed as shown in FIG. 4(a)."}
{"text": "1. Field of the Invention\nThe present invention relates to social proxies, and more particularly towards a marketplace social proxy for the facilitation of business-to-consumer interaction.\n2. Description of Prior Art\nThe growth of on-line consumer support and sales has led to important questions regarding consumer education, attitude change, persuasion, decision making, and satisfaction in on-line environments. Current on-line sales environments isolate consumers from one another. These sales environments provide only basic sales service, for example, offering a picture of the product/service and a general description. Some independent services (e.g., cnet.com) provide consumer-to-consumer interaction via bulletin boards to serve as distributed knowledge (opinion) reservoirs. However, no known direct sales approach has implemented a system or method for consumer-to-consumer interaction.\nIn the area of online user interaction, there are several efforts working towards creating a computer interface capable of displaying intangible qualities typically associated with live interaction. In particular, FTP, NNTP, IRC (e.g., The Palace), Babble, Loom, Coordinator, usage summaries for web sites, Ebay and Amazon.com's user interface. These systems include elements which relate to attributes normally associated with a physical interaction. For example, by posting the number of visitors to a web site (e.g., usage summaries), a user can have a feeling of belonging to a larger group of people interested in that web site. These and other systems are described below.\nFTP allows online users “anonymous”—to have access to a collection of documents; one or more users being authorized to add documents to the collection. (See Ed Krol, The Whole Internet Catalog. O'Reilly and Associates, Inc. Sebastopol, Calif. 1992.) Interfaces to FTP-accessible collections of documents (e.g., bulletin boards) do not provide dynamically updated, graphical representations of the activity of a given online marketplace environment.\nNNTP or network news, provides a venue in which users can asynchronously post messages and responses into administered news groups (i.e., predefined groups of messages which are meant to relate to specified topic, e.g., all articles in alt.sport.soccer should deal with soccer), every post indicating the ID of the contributor. Lotus Notes' discussion databases provide a functionally equivalent form of online interaction. Here again, no resource based on these technologies provide business users with a dynamically updated, graphical representation of the activity of a given online business environment.\nInternet Relay Chat (IRC) is similar to NNTP in that it allows users to post messages and responses in predefined topic areas—called channels in IRC. The main difference with IRC is that the interactions are synchronous rather than asynchronous as in NNTP.\nThe Palace, is a collection of Internet chat rooms. The site uses a graphical display and user icons. A Palace Site is an online community where users can come together and chat while being represented by a graphical image. There are thousands of existing Palace Sites. Users can connect to The Palace Sites with The Palace Viewer (PTV) or The Palace User Software. Currently, channels available include: Welcome, TV, Movies, Music, Romance, Teens, and International. Within each channel is a collection of related Palace Sites. For example, within the Music channel, specific Palaace Sites may exist for particular genres or artists. TPV allows users to visit Palace Sites through their web browsers. With TPV and The Palace Server, a user can add a Palace community to a web site. The Palace User Software enables a user to connect to all the different Palace Sites located in Palace Space.\nOther implemented forms providing synchronous online communication also exist, for example, Instant Messaging from Apple, and the discussion section of Groove from Groove Networks, Inc. Although these applications and systems provide user with a list of all operational clients, none provides a dynamically updated graphical representation indicating the type or level of activity of the users.\nLoom, a visualization tool for Usenet (NNTP) news groups, provides a technique for displaying the emotional mood (e.g., hostile, happy) of an NNTP-based online community (e.g., a Usenet news group) by analyzing the content of its interactions. (See Judith Donath et al. “Visualizing Conversation” published in the Journal of Computer Mediated Conversation. Volume 4, number 4, June 1999.) Although this utility can analyze online interactions and provide a graphical representation of aspects of the interactions, it only provides this information asynchronously, specifically by analyzing the overall content records from each NNTP room. Thus, the graphical representation it provides is not dynamically updated.\nBabble provides dynamic indicators of the presence and activity of all operational users with respect to the available topics (i.e., discussion areas created by the users). These indicators are computed using the activities of the participants (e.g., connections, postings, and topic switches). (See Thomas Erickson, et al. “Socially Translucent Systems: Social Proxies, Persistent Conversation, and the Design of ‘Babble.’” published in Human Factors in Computing Systems: The Proceedings of CHI '99. ACM Press, 1999.) Although lightly structured styles of interaction are frequently adopted (e.g., interaction in the “-Commons Area-” is informal conversation, while interaction in the “Babble Problems” topic consists of serious question/answer dialogs), no way is provided to customize the graphical representation or the information it renders so that the specific critical qualities and quantities of an online marketplace are (dynamically) displayed.\nOnline games, like chess and bridge (e.g., Chessmaster 6000 by Mindscape, Inc.), provide structured and enforced styles of interaction, but not ones that dynamically and graphically indicate the activity of a given online marketplaces.\nCoordinator, a method of structuring human communications, provides an electronic mail-based system that allows users to structure conversations and track tasks. For example, a typical interaction begins with a “Request” message from person A to person B, requesting something from person B by a certain date. This e-mail note asks Person B to respond with a “Promise” message (promising to perform the action), with a “Decline” message (declining to perform the action), or with a “Counteroffer” message (offering to perform the action by a different date or to perform a different action). If B promises to do the action, then a typical conversation might continue with B, eventually sending a “Report completion” message (indicating that the action has been performed) and A replying with a “Close” message (indicating that the action was performed satisfactorily). (See: Flores et al. “Method and Apparatus for Structuring and Managing Human Communications By Explicitly Defining the Types of Communications Permitted Between Participants.” U.S. Pat. No. 5,208,748, May 4, 1993.)\nThe Oval implementation of Coordinator extends the base functionality by allowing end-users to modify interaction rules mid interaction (see Malone et al. “Experiments with Oval: A Radically Tailorable Tool for Cooperative Work” via URL http://ccs.mit.edu/papers/CCSWP181/index.html#4b).\nThus, Coordinator, implemented with Oval allows users to define and maintain structured styles of online interaction. These two facilities still fail to provide a dynamically updated, graphical representation of the state and activity within a given online marketplace; and handle interactions where messages are posted to a group of recipients, all of whom are not known a priori (e.g., as is the case in many online business environment, i.e., one often does not know who one's customers will be at the start of the day).\nIn the online auctions provided by Ebay.com., users can get data concerning the activity of other buyers and sellers (e.g., how many transactions they have successfully completed). Although this provides an indication of user activity in a online marketplace, no dynamically updated, graphical representation is provided.\nAmazon.com also provides information concerning the activities and comments from it customers. E.g., for any given product sold by Amazon.com, one can learn what other buyers who have purchased the given product have also bought (e.g., customers who purchased the first Harry Potter book also purchased every other book about Harry Potter). Although this provides indication of user activity, Amazon.com does not provide any graphical representation, nor does it provide information that is dynamically updated.\nTherefore, a need exists for a method and system for implementing a graphical online environment including consumer and business proxies for providing a dynamically updated, graphic view (representation) of current marketplace conditions."}
{"text": "1. Field of the Invention\nThe present invention relates generally to DC to DC converters and specifically to a circuit to recover energies in snubbing capacitors by using auxiliary switch(es) in a single-ended forward converter. More particularly, it relates to such a circuit in single-ended interleaved or non-interleaved forward converters. Most especially, it relates to such a circuit which clamps the peak primary switch voltages to below a predetermined voltage, independent of output power, and properly resets the power transformers, introducing virtually no losses.\n2. Description of the Prior Art\nThe demand for distributed power systems in computer and communications industries has mandated the power conversion from a low dc bus voltage, typically 48 volts, to voltage levels suitable for integrated circuits, such as 5.1 volts and 3.3 volts. In the design of such DC-DC converters, one of the safety requirements for meeting CLASS 3 per IEC Publication 950 is that the peak primary voltage cannot exceed 140 volts between any two points under any operating condition.\nA single-ended forward-type pulse width modulated (PWM) converter is a prime candidate for this application. At higher power levels, interleaved type forward converters are used. The converter usually operates with an input voltage range of 40-60 volts for accommodating a possible battery back up. To prevent the peak primary voltage from exceeding 140 volts, a large snubbing capacitance has to be used across the main switch. The energy in the snubbing capacitance is lost when the main switch turns on, and thus the converter's efficiency suffers. As the output power increases, the situation becomes more severe because the energy in the leakage inductance increases and as a result the snubbing capacitance has to be larger, resulting in greater loss of efficiency as the larger snubbing capacitance is discharged through the primary switch.\nA distributed power system usually consists of a front-end power stage, which converts typical ac outlet voltage to a safe low dc voltage and performs power factor correction to the ac line. The low dc voltage is then distributed to the loads, where small local converters are used to convert the low dc voltage to proper levels, for example for use with integrated circuits. Such an architecture reduces the hazards of high voltage distribution, provides built-in redundancy over systems in which voltages are centrally generated and distributed, and has better dynamic responses to the loads. The most popular dc voltage selected for low dc voltage distribution systems is 48 volts. Many converter modules converting 48 volts to proper voltage levels for integrated circuits have recently been developed and are commercially available. Various converter topologies have been adopted in the converter modules. The single-ended forward topology with resonant reset, as shown in FIG. 1, is among the most popular. For higher power levels, interleaved single-ended forward converters, as shown in FIG. 2, are used. FIGS. 1A and 2A illustrate some of the waveforms present in the circuits of FIGS. 1 and 2, respectively.\nWhen the main switch of a single-ended forward converter turns off, energies stored in the magnetizing and leakage inductances tend to resonate between the inductances and the output capacitance of the main switch, such as a MOSFET, and generate voltage spikes and high-frequency ringing. Note that a MOSFET effectively has a capacitor C.sub.M between its drain and source terminals, and a diode D.sub.M from its source to drain terminal as illustrated by dotted lines in conjunction with a MOSFET in the insert in FIG. 1. These elements are not explicitly illustrated in the circuit schematics, but are understood to be present by those skilled in the art.\nTraditionally, resistance-capacitance (R-C) or resistance-diode-capacitance (R-D-C) snubber circuits are used to provide a low impedance path to limit the spikes caused by high change of current with respect to time (di/dt), and/or to damp the ringing. An R-C snubber circuit is a resistor connected to the MOSFET main switch drain in series with a capacitor connected to the source. An R-D-C snubber circuit has the structure of the R-C snubber circuit but with an additional diode in parallel with the resistor and with its anode connected to the drain and the cathode connected to the capacitor. These types of snubber circuits, however, are very lossy and inefficient.\nAnother traditional approach is to add only external capacitance across the main switch to decrease the characteristic impedance and the resonant frequency. This results in a reduction of the spikes and the ringing by significantly reducing the resonant frequency. However, the introduction of the additional capacitance disadvantageously increases the turn-on energy losses of the main switches (the capacitor energies are dissipated when the main switch turns on), and the converter's efficiency suffers. To meet CLASS 3 per IEC Publication 950 safety requirements, large capacitances are needed to prevent the peak main switch voltage(s) from exceeding 140 V under all operating conditions. The requirement for large capacitances, typically about 20 to 50 times the output capacitance of the MOSFET, serving as the main switch, depending on the energy of the transformer leakage induction, makes this approach unattractive. Although the spikes and ringing can also be reduced by increasing the magnetizing inductance, there is also a practical limitation, as well as cost concerns, to making large magnetizing inductances.\nAnother approach, the zero-voltage switched-snubber concept, may effectively clamp the voltage across the main switch or the output rectifier diodes in PWM converters to a specified level while recovering the energies in the snubbing capacitor. The zero-voltage switching concept requires that the switching of the main switch occur when the voltage of the snubbing capacitor is zero or a very low voltage. As the voltage across the main switch at the time of its switching on increases, the required snubbing capacitance increases as the square of the voltage. A larger snubbing capacitance means greater energy loss when switching. The zero-voltage switched snubber concept allows the use of lower voltage rated switching devices as the main switch, thus reducing the devices' conduction losses and improving the overall efficiency of the converter. However, although zero-voltage-switching can be achieved using switched snubbers, such implementation requires additional gate-drive logic and extra drivers for the auxiliary switch or switches. This additional gate-drive logic and driver(s) is not commercially available as a single integrated device, and thus must be made specifically for this purpose, resulting generally in greater size and cost. Furthermore, the transformers having specially designed characteristics may be needed to operate in conjunction with these additional components.\nThe typical circuit operation of a single-ended forward converter as in FIG. 1 is now described. When MOSFET main switch M1 turns on, capacitor CS1 is totally discharged (V.sub.CS1 =0), diode DR1 is conducting, and diode DR2 is off. The majority of the input voltage, E, given by the expression ##EQU1## where: LM1=transformer magnitizing inductance,\nLK1=transformer leakage inductance, PA1 E=voltage provided by voltage source E, PA1 N=the primary to secondary transformer turns ratio PA1 E=voltage provided by voltage source E, PA1 V.sub.DR2 =voltage across diode DR2 PA1 i.sub.p =reflected load current in the transformer primary winding, PA1 I.sub.0 =load current in RL, and PA1 N=the primary to secondary transformer turns ratio.\nis transformed to the transformer secondary, and the voltage across diode DR2 is given by, EQU V.sub.DR2 .perspectiveto.E/N (2)\nwhere:\nand the load current is reflected back to the transformer primary. Therefore, the reflected load current in the primary, i.sub.p, is related to the load current, I.sub.0, by the expression, EQU i.sub.p .perspectiveto.I.sub.0 /N (3)\nwhere:\nWhen main switch M1 is turned off, the current in leakage inductor LK1, designated i.sub.LK1 (which includes the magnetizing current i.sub.LM1, and the reflected load current, i.sub.p) is diverted to capacitor CS1. The voltage V.sub.CS1 across capacitor CS1 rises, and the voltage transformed to the secondary decreases. When V.sub.CS1 exceeds the magnitude of E, negative current flows in the primary inducing voltage in the secondary such that diode DR2 is turned on and diode DR1 starts to be commutated. During the commutation interval, transformer T1 secondary winding is effectively shorted because both diode DR1 and DR2 are conducting (for a very brief time) and inductor current i.sub.LK1 resonates with capacitor voltage V.sub.CS1.\nWhen diode DR1 turns off, output current freewheels through diode DR2 and the capacitor voltage V.sub.CS1 begins to resonate with i.sub.LM1. (At this instant, the magnitizing current i.sub.LM1 is usually negative. In this event, capacitor voltage V.sub.CS1 decreases.) When V.sub.CS1 decreases below E, positive current flows in the primary, inducing a voltage in the secondary and causing diode DR1 to conduct and commutate diode DR2. However, when diode DR1 starts conducting, reflected current i.sub.P forces V.sub.CS1 to increase and try to commutate diode DR1. If diode DR1 is commutated, the negative magnitizing current i.sub.LM1 forces V.sub.CS1 to decrease again. This mechanism forces V.sub.CS1 to resonate with i.sub.LK1 around the input voltage level until MOSFET main switch M1 is turned on (triggered) by signal V.sub.G1 applied to its gate. At that instant, V.sub.CS1 decreases to zero and the energy in capacitor CS1 is dissipated through MOSFET main switch M1. Diode DR1 is forced to conduct and diode DR2 is commutated. Load current I.sub.0 is again reflected to the primary side of transformer T1 until MOSFET main switch M1 is turned off by trigger signal V.sub.G1, which starts another cycle of the operation."}
{"text": "This invention relates to an improved secondary purification process for phosphoric acid. Specifically, it relates to such a process where a portion of the filtrate from the calcium sulfate dihydrate filter is recycled to the same filter.\nSuperphosphoric acid (SPA), which is used in the manufacture of fertilizers, is phosphoric acid containing about 68 to about 70% by weight of P.sub.2 O.sub.5. It can be manufactured by dissolving calcium phosphate rock with sulfuric acid, producing a slurry of calcium sulfate hemihydrate in a solution of phosphoric acid. The slurry is filtered and the calcium sulfate hemihydrate filter cake is discarded. The filtered acid is evaporated to about 48% P.sub.2 O.sub.5 for further purification.\nSome phosphate rock used in this process contains a significant quantity of magnesium, which is dissolved by the sulfuric acid. If the magnesium is left in solution with the phosphoric acid, magnesium pyrophosphate forms when the solution is evaporated to SPA. When the SPA is ammoniated, an ammonium magnesium pyrophosphate sludge forms which creates a serious pluggage problem for liquid fertilizer manufacturers.\nIn order to avoid this problem, the magnesium is removed from the phosphoric acid. This can be accomplished by the addition of synspar, essentially a calcium fluoride compound made as described in U.S. Pat. Nos. 4,171,342 and 4,243,643, herein incorporated by reference. The addition of the synspar to the phosphoric acid slurry results in the precipitation of Ralstonite, MgAlF.sub.5 (aluminum is normally present in the rock and is dissolved by the sulfuric acid), and hemihydrate gypsum.\nTo remove the Ralstonite and hemihydrate gypsum, the slurry is filtered on a rotary drum vacuum filter. The hemihydrate cake containing Ralstonite is discharged to a hydration tank, and is reslurried with pond water. The hemihydrate gypsum dissolves and reforms as dihydrate gypsum in this tank. The slurry is then filtered on a second rotary drum vacuum filter. The filtrate is mixed with the filtrate from the first drum filter for further processing. The dihydrate cake containing Ralstonite is disposed of.\nIn order to maximize the P.sub.2 O.sub.5 recovery, it was, until now, necessary to maintain a minimum of 2% free sulfate ion in the hydration tank. As long as the free sulfate ion concentration in the hydration tank was maintained at least 2%, recrystallization of the calcium sulfate hemihydrate as calcium sulfate dihydrate occurred even though the residence time of the slurry in the hydration tank often varied considerably. But when the free sulfate ion concentration fell below about 2% more time was required to convert the hemihydrate crystals to dihydrate crystals, and frequently the residence time in the hydration tank was not long enough for a complete conversion to occur. When that happened calcium sulfate hemihydrate entered the calcium sulfate dihydrate filter and formed calcium sulfate dihydrate on the filter, plugging the filter.\nIn order to maintain a free sulfate ion concentration of at least 2% in the hydration tank, it was necessary to add sulfate ion (as sulfuric acid) in excess of the amount required to precipitate the dissolved calcium as calcium sulfate. However, if calcium phosphate rock having a high concentration of magnesium was used in the process, it was necessary to add extra synspar to remove the extra magnesium, and, since the calcium in the synspar removes sulfate ion as calcium sulfate, the free sulfate ion concentration would sometimes fall to less than 2%, or even less than 1%. When the free sulfate ion concentration was less than 2%, the residence time of the calcium sulfate hemihydrate slurry in the hydration tank was frequently not long enough to convert all of the calcium sulfate hemihydrate to calcium sulfate dihydrate. The unconverted calcium sulfate hemihydrate would then enter the calcium sulfate dihydrate filter where it would form the calcium sulfate dihydrate and clog the filter. If the residence time in the hydration tank was increased by shutting off the flow of the slurry from the hydration tank to the calcium sulfate dihydrate filter, the calcium sulfate dihydrate filter would become \"dry\" and would have to be shut down and restarted again when the required residence time had passed and slurry was again available for filtering. This resulted in P.sub.2 O.sub.5 losses in the lines and equipment, known as \"tramp\" losses. There did not appear to be any way of solving this problem without reducing the efficiency of the process and lowering the recovery of P.sub.2 O.sub.5."}
{"text": "1. Field of the Disclosure\nThis disclosure relates generally to apparatus and methods for dispensing a viscous material on a substrate, such as a printed circuit board, and more particularly to a system and related methods for accurately dispensing material on a circuit board based on detecting an edge of the circuit board.\n2. Discussion of Related Art\nThere are several types of prior art dispensing systems used for dispensing precise amounts of liquid or paste for a variety of applications. One such application is the assembly of integrated circuit chips and other electronic components onto circuit board substrates. In this application, automated dispensing systems are used for dispensing very small amounts, or dots, of viscous material onto a circuit board. The viscous material may include liquid epoxy or solder paste, or some other related material. Prior to performing a dispensing operation, the circuit board must be aligned or otherwise registered with a dispenser of the dispensing system. In one known method, this may be achieved by employing a vision system of the dispensing system to verify the location of landmarks, otherwise known as fiducials, on the circuit board. Specifically, to align the circuit board with the dispensing unit of the dispensing system, images of at least two fiducials are taken by a camera of the vision system. If the circuit board is out of position, a gantry capable of moving the dispenser may be manipulated to account for the actual position of the circuit board. In another embodiment, a support surface upon which the circuit board rests may be manipulated to accurately position the circuit board prior to performing a dispense operation.\nOne issue associated with the use of identifying fiducials to align the circuit board with the dispensing unit is when the area requiring dispensing is not well defined with reference to the fiducials. For example, if the area to dispense is not referenced with respect to the fiducials, the resulting dispense operation may be very inaccurate, and under many instances, may not meet the processing requirements."}
{"text": "The embodiments disclosed herein are drawn generally to a dual-sensor stethoscope that promotes anti-sepsis, stereoscopy, and advanced sound processing. In particular, a stethoscope with dual-sensors on a T-shaped support is disclosed, wherein at least one sensor is electronic.\nExisting stethoscopes are currently utilized to auscultate or listen to physiologic sounds within the body. Auscultation with existing stethoscopes is currently performed by intermittently applying a stethoscope to the body surface through which the clinician hears various sounds. Intermittent auscultation may be thought of as a relatively benign procedure. However, several disadvantages and hazards are associated with the use of existing stethoscopes. First, patients undergoing surgery may have the sterile field invaded thereby risking infection in order for the clinician to auscultate the chest. To avoid cross-contamination between patients, many clinicians are forced to carry multiple stethoscopes. Additionally, even when a sterile stethoscope is used, it can transfer pathogens from a first location on a patient's body to a second location during typical auscultation. Furthermore, even in non-surgical environments, transmission of the cold virus is primarily through touch. A clinician's hand can touch the head of a stethoscope, which then touches a patient and vice versa so as to spread the virus. These issues exist whether the stethoscope uses an acoustic (mechanical) sensor or an electronic sensor.\nAnother disadvantage of known stethoscopes is that patients are frequently awakened and disturbed so that the clinician may apply a cold stethoscope to the patient's chest to monitor vital signs. Studies have shown serious developmental abnormalities in newborn infants who are frequently disturbed to auscultate heart and lung sounds with known stethoscopes. Another disadvantage of existing stethoscopes is that the quality of sound wave transmission is dependent upon an airtight seal between the stethoscope and the skin, typically requiring the clinician to touch, and possibly contaminate, the sensor. In the absence of an airtight seal, background noise is inadvertently detected and physiologic sound transmission is impaired. Finally, another disadvantage of existing stethoscopes is that most are not capable of generating positive or constructive interference, filtering certain frequencies, or providing other processing of physiologic sound waves."}
{"text": "Flash memory has gained wide acceptance for its non-volatile storage, which is ideal for portable devices that may lose power, since the data is not lost when stored in the flash memory. Flash memories are constructed from electrically-erasable programmable read-only memory (EEPROM) cells.\nRather than use a randomly-addressable scheme such as is common with dynamic-random-access memory (DRAM), many flash memories use a block-based addressing where a command and an address are sent over the data bus and then a block of data is read or written. Since the data bus is also used to send commands and addresses, fewer pins are needed on the flash-memory chip, reducing cost. Thus flash memory is often used as a mass-storage device rather than a randomly-addressable device.\nUniversal-Serial-Bus (USB) has become a popular standard interface for connecting peripherals to a host such as a personal computer (PC). USB-based flash-memory storage devices or “drives” have been developed to transport data from one host to another, replacing floppy disks. While large external flash drives may be used, smaller USB flash drives known as key-chain or key drives have been a rapidly growing market.\nA USB flash-memory device can be constructed from a microcontroller, a flash-memory controller or interface, and one or more flash-memory chips. A serial interface on the microcontroller connects to the USB bus to the host, and data from the serial interface is transferred through the microcontroller to the flash controller and then written to the flash-memory chips.\nThe microcontroller usually contains an internal ROM with a control program that is read by the internal central processing unit (CPU) of the microcontroller when the microcontroller is booted or powered up. Once initialized with the control program, the CPU can control data transfers between the serial interface and the flash controller.\nSometimes the user may desire to connect to more than one USB flash-memory device. The user can install a USB hub, and then plug the USB flash-memory devices into the USB hub's downstream ports. USB hubs allow one USB port on a host to fan out to multiple end USB devices or endpoints. A basic USB hub has a repeater that repeats data from the host to all down-stream devices, while more intelligent hubs based on the USB 2.0 standard can buffer data to different down-stream ports.\nThe parent application, now U.S. Pat. No. 7,103,684, disclosed a USB flash drive that did not need a read-only memory (ROM) for booting. The microcontroller in the USB flash drive was able to read boot code from the block-addressable flash memory and transfer the boot code to the microcontroller's RAM for execution.\nAnother popular bus standard is Multi-Media Card (MMC). An extension of MMC is known as Secure Digital (SD). MMC and SD flash devices are common today. It is desired to extend the ROM-less flash-memory-drive microcontroller of the parent application to MMC, SD, Memory Stick (MS) and other similar portable buses.\nFIG. 1 shows a prior-art MMC flash controller that connects to multiple flash-memory devices. Host 10 includes MMC host controller 12 that generates transactions to MMC devices over MMC bus 18 using the MMC protocol. MMC controller 20 is connected to a cable containing MMC bus 18. MMC controller 20 fans out commands on MMC bus 18 to several downstream flash-memories that connect over additional bus segments.\nThree flash-memory systems 14, 15, 16 are connected to MMC controller 20 by bus segments. Flash-memory system 14 can be accessed by MMC host controller 12 through MMC controller 20. Since MMC controller 20 passes all host transfers through to downstream devices, flash-memory system 15 is readable to host 10 as a second flash drive, while flash-memory system 16 is readable to the host as a third flash drive. Alternately, all three flash-memory systems 14, 15, 16 may appear to host 10 as a single memory.\nSome board or device manufacturers may integrate MMC controller 20 together with flash-memory systems 14, 15, 16 on a single MMC flash card or box. However, this can be expensive when MMC flash-memory systems 14, 15, 16 are flash-memory chips, since each chip may have many pins. For example, a flash-memory chip with an 8-bit or 16-bit data bus may have 48 total pins. This can increase the size of the MMC flash device. Power consumption is higher due to the large number of data lines in the parallel buses to each flash-memory chip.\nWhat is desired is to integrate a microcontroller with a flash-memory array. It is desired to have a wide internal bus from the microcontroller to the flash-memory array to improve the data bandwidth while having few external pins to reduce cost and required board space.\nIt is further desired to eliminate the internal ROM on the microcontroller. Instead of booting from the internal ROM, it is desired to use a control program stored in the flash-memory array. However, it is also desired to use a block-addressed rather than a randomly-addressable array for the flash storage."}
{"text": "1. Field of the Invention\nThe present invention relates to a computer peripheral, and more particularly to a quick-assembly computer peripheral that has a bottom casing and a top casing quickly attached to the bottom casing to facilitate a semi-automated or fully-automated manufacturing process, thus lowering manufacturing costs and raising the production rate of the computer peripheral.\n2. Description of Related Art\nComputer peripherals comprise keyboards, mice, trackballs, pen tablets and the like. A conventional mouse has a bottom casing, a top casing, a printed circuit board (PCB) and a cable. The top and bottom casings are combined by screws and form an inner chamber. The PCB is mounted in the inner chamber and has multiple preformed mounting holes. The cable has multiple cores soldered respectively in the mounting holes to complete the electrical connection. However, mounting the top and bottom casings by screws and individually soldering the cores to the PCB wastes time, increases manufacturing costs and lowers production rate of the mouse.\nTo overcome the shortcomings, the present invention provides a quick-assembly computer peripheral to mitigate or obviate the aforementioned problems."}
{"text": "For example, a synchronous type adaptive filter which cancels from a main input a periodic noise originating from a revolving drum motor of a camera-integrated type video tape recorder (hereinafter abbreviated to VTR) is proposed in Patent Gazette of Japanese Published Patent Application No. H11-176113 and so on. In the adaptive filter employed in a technology disclosed in this gazette, the noise is reduced by renewing and converging filter coefficients adaptively to the target noise occurring at a constant period (for example, 150 Hz) along with the above-described motor revolutions.\nFIG. 6 shows a block diagram of a conventional adaptive noise reduction circuit, an outline of which will be described. In FIG. 6, an input terminal 1 receives a main audio input S whose noise is to be reduced. In actual equipment, for example, noise N schematically shown is supplied to an input terminal 2 and is simply added to the main audio input S at an adder 8, which causes the main audio input S to be mixed with the noise N. The adaptive noise reduction method reduces the noise N from the main audio input S mixed with the noise N, and therefore, the main audio input S mixed with the noise N is supplied to a (+) side terminal of an adder 9.\nOn the other hand, a reference input X highly correlated with the noise N is supplied to an input terminal 3. The reference input X is supplied to an adaptive filter 6, where an adaptive filter output Y that is approximate to the noise N is formed by adaptation processing. Then, an adaptive filter output Y thus formed is supplied to a (−) side terminal of the adder 9 and is subtracted there from the main audio input S mixed with the noise N. Accordingly, an audio output S^ is derived from the adder 9, from which the noise N is removed by the adaptive filter output Y as is shown by the following expression (0).S^=S+N−Y   expression (0)\nThis audio output S^ is an audio signal of originally aimed, from which the noise N is removed and is output to an output terminal 10. Simultaneously, the audio output S^ (a residual signal E) is fed back through a step gain 7 to be used in adaptation processing. Specifically, the residual signal E is supplied to, for example, a least mean square (hereinafter abbreviated to LMS) operation processing circuit 5 together with the reference input X, where an operation of coefficients of the adaptive filter 6 is performed so that, for example, noise power of the residual signal E may become minimum.\nFurther, the above-described adaptive filter 6 will be described below in detail with reference to a block diagram of FIG. 7. Note that, while various methods have been proposed for algorithm of the adaptive filter, the above-described LMS method is often employed because its convergence speed is comparatively fast and a scale of operational circuit is small in general. Moreover, those circuits can be processed with hardware whose circuit is all made up of digital signal processor (hereinafter abbreviated to DPS), digital large scale integrated circuit (hereinafter abbreviated to digital LSI) and the like, or software using microcomputer.\nHereupon, in FIG. 7, the reference input X is supplied to the adaptive filter 6 surrounded by a broken line as well as the LMS operation processor circuit 5. The adaptive filter 6 is conventionally includes a FIR digital filter which has about several hundred taps, and the filter coefficient W in each of the taps is adaptively renewed in accordance with, for example, the LMS algorithm. In FIG. 7, the FIR filter having (m+1) taps is shown, which is provided with unit delay means 111 to 11m each having a delay Z−1 of unit sampling time.\nTherefore, delayed signals X0 to Xm are obtained from these unit delay means 111 to 11m, respectively. These signals X0 to Xm are supplied to multipliers 120 to 12m for multiplying them by coefficients. Moreover, to these multipliers 120 to 12m are supplied the adaptive filter coefficients W0 to Wm formed in, for example, the LMS operation processing circuit 5. All outputs of these multipliers 120 to 12m are added in an adder 13 to be output as the adaptive filter output Y.\nThe adaptive filter output Y is represented by the following expression (1).\n                    Y        =                              ∑                          j              =              0                        m                    ⁢                      (                          Wj              ·              XJ                        )                                              expression        ⁢                                  ⁢                  (          1          )                    \nFurther, in the LMS operation processing each of adaptive filter coefficients W0 to Wm is renewed in accordance with the following expression (2) from the above-described reference input X and residual signal E.Wk+1=Wk+2μ•Ek•Xk   expression (2)\nIn this expression (2), the small letter k represents passing time, as an example, assuming that a value k is renewed by every unit sampling, the value k indicates the kth sampling and a value (k+1) indicates the (k+1)th sampling.\nThe value μ is a coefficient given by the step gain 7. The value μ is called a step gain or step size which is a parameter that determines a convergence speed in the LMS algorithm. It is noted that if the value μ is large, the convergence becomes fast, however accuracy after convergence falls; inversely, if the value μ is small, the convergence becomes slow, however accuracy after convergence rises. For this reason, the value μ is set at an optimum according to an adaptive system condition in use and the like.\nIn this way, according to the above-described apparatus, in the LMS operation processing the adaptive filter coefficient W in the adaptive filter is renewed according to the expression (2) such that the signal highly correlated with the reference input X and included in the residual signal E is kept to a minimum, so that a component of noise N contained in the main audio input S can be kept to a minimum by inputting to the reference input X the above signal correlated with noise N. In other words, between the LMS operation processor circuit 5 and adaptive filter 6, such a feedback loop that makes minimum of a component of the residual signal E is formed.\nA conventional fixed-period noise reduction block will further be described with reference to FIG. 8. In the following description, elements of the same function as in FIG. 6 are denoted by the same reference numerals.\nFirst, similarly to FIG. 6, the main audio input S whose noise is to be reduced is supplied to an input terminal 1. In actual equipment, for example, noise N is schematically supplied to the input terminal 2 and is added simply to the main audio input S in the adder 8 to make the main audio input S mixed with noise N. A pseudo-noise signal Yk is supplied to a (−) side terminal of the adder 9 from an adaptive signal processor 20 surrounded by a broken line, and by subtracting the pseudo-noise signal Yk, only the main audio input S is obtained at the output terminal 10.\nAt the same time, the above-described error signal through the step gain 7 is supplied to the adaptive signal processor 20. Further, the fixed-period pulse signal as the reference input signal from an input terminal 3 and a sampling clock from an input terminal 21 are supplied to the adaptive signal processor 20. Hereupon, the adaptive signal processor 20 includes the same adaptive filter as in FIG. 7, which renews the coefficients in accordance with the LMS algorithm, and processing thereof is performed in synchronism with the sampling clock from the input terminal 21.\nIn addition, the sampling clock corresponds with a sampling frequency of the main audio input S and noise N. In FIG. 8, the sampling clock from the input terminal 21 and the fixed-period pulse signal from the input terminal 3 are inputted to a counter 22 within the adaptive signal processor 20, and a period of the input is counted by the sampling clock and the counted value thereof is supplied to a timing generator 24 to generate a predetermined timing pulse.\nThen, from the timing pulse, an Xk address is generated in turn from 0 to m in a read-address generator 23 and Xk−1 address is generated in turn from 0 to m in a write-address generator 25. These Xk addresses and Xk−1 addresses are respectively input as the read address and write address to an accumulator 26 composed of static RAM (hereinafter abbreviated to SRAM) or the like.\nFurther, the accumulator 26 has registers of, for example, m+1 words (m+1 taps) at the maximum, each of which has a predetermined bit length. Then, the accumulator 26 of m+1 taps is designed that an adaptive coefficient W is read from or written into the respective specified addresses with the predetermined timing within one period, according to the read address Xk or the write address Xk−1.\nFurthermore, to one terminal of an adder 28 is supplied 2 μEk which is obtained by multiplying the above-described error signal Ek by the step-gain p, and to the other terminal of the adder 28 is supplied data Wk read from the address Xk of accumulator 26. An output of the adder 28 obtained by adding both the inputs is then delayed by a unit sampling time 27, and the resulted signal is written into the Xk−1 address of the accumulator 26. Likewise, the pseudo-noise signal Yk of one period before is read from the Xk address.\nAccordingly, in FIG. 8, the adaptive filter coefficient W is renewed using the above-described expression (1) such that the noise component which is contained in the error signal E and is highly correlated with the drum reference signal X may always become minimum, and a main signal output whose noise is always reduced is obtained at the output terminal 10.\nNext, an example of addressing by the fixed-period adaptive filter shown in FIG. 8 will be described with reference to FIG. 9. In this case, 0 to m words of the accumulator 26 in FIG. 8 is first formed in the shape of a ring. In other words, an address next to Adr:m is made to be Adr:0, or an address preceding to Adr:0 is made to be Adr:m. Assuming that a time of one period is T[s] and a sampling frequency is S[Hz] with respect to the number of sampling clock within one fixed period, the value m has the following relation.m≈S•T   expression (3)\nThen, the write-address generator 25 controls data within the ring-shaped memory to move in the direction shown with an arrow by one address at every sampling clock, and to fix relatively a read-address position and a write-address position at positions shown in the figure. Accordingly, a write signal as shown in FIG. 10A, for example, and a read signal delayed by a predetermined amount as shown in FIG. 10B are obtained. A pitch of both the signals is definite and becomes a fixed period m.\nSpecifically, according to the above-described apparatus, the noise can be reduced by renewing and converging the filter coefficients adaptively to target noise occurring at a constant period (for example, 150 Hz) as a rotary drum motor of, for example, a camera-integrated VTR revolves.\nHowever, in case of the variable period noise in which a period of, for example, target noise varies, it is necessary to renew the number of taps and coefficients of the above-described adaptive filter in accordance with the periodic variation. Specifically, with respect to the target noise occurring at a definite period, the noise is reduced with such adaptive filter that has, for example, as many taps as a quotient got when that period divided by the sampling period, whereas if the period varies, the number of taps must be changed. Moreover, when the period varies, it is assumed that a waveform of the target noise will also change.\nTo cope with this problem, it is conceived that the adaptive filter coefficients must be renewed corresponding to the periodic variation. In general, however, because the step-size (or step-gain) that is a parameter determining a time required for renewal of the filter coefficient is determined taking into account the influence from external disturbance and the like, it is impossible to increase the step-size (or step-gain) corresponding to the periodic variation. For this reason, there is a possibility that, when the period varies fast for example, there may happen a case where the adaptive filter cannot follow, thus making the noise reduction impossible disadvantageously.\nAdditionally, an example of the variable period noise is an electromagnetic noise, noise, and a vibration noise originating from the disc motor in DVD-RAM. In this case, the disc motor adopts a zoned constant linear velocity (hereinafter abbreviated to ZCLV) system for a revolution control system. This system divides a recording area into zones according to a radial position of the disc and sets the number of revolution of disc such that a recording density in each zone may become almost constant, thus causing the number of revolution to vary in each zone.\nIn this example, the number of revolution varies in the range of, for example, 3246 (at inner circumference) rpm to 1375 (at outer circumference) rpm. Thus, during the movement between zones at seek time for example, the number of revolution must be changed rapidly, which causes the above-described variable period noise to occur depending on a change in the number of revolution. In addition, even though the revolution control system is the conventional CLV system, the number of revolution is rapidly varied at the seek time or the like, so the variable period noise will occur depending on a change in the number of revolution.\nThis application has been made in view of those points and aims to solve a problem in which, with the conventional adaptive noise reduction method and apparatus, it is difficult to make the adaptive filter coefficients rapidly follow the variable period noise whose period changes rapidly, thus making it impossible to reduce efficiently the variable period noise from the main input."}
{"text": "Ocular surgery, whether reconstructive, cosmetic, palliative, or otherwise, is highly patient specific. Even though most surgery patients have the same basic ocular architecture, every eye has its own set of specific features and dimensions that in certain cases may be significantly different from those of expected norms. As a result, surgeons must rely upon their individual experience and skills to adapt whatever surgical techniques they are practicing to the individual requirements as determined by each patient's unique ocular structural features and dimensions.\nTo date, this individualized surgical adaptation has been accomplished essentially through freehand and best guess techniques based upon a pre-surgery examination and evaluation of the individual patient's ocular region. This examination may include preliminary measurements as well as the surgeon making reference markings directly on the patient's ocular tissues with a pen or other form of dye or ink marking. Then, after the patient has been prepared and placed in position for surgery, typically in a supine or prone position as opposed to the often vertical positioning of the patient during the pre-surgery examinations, the surgeon adapts the placement and configuration of the initial surgical incisions to the actual physical dimensions and circumstances found in the patient as the surgical procedure progresses. As a result, many initial measurements or reference markings on the patient's ocular tissues are at best a general guide as to where to begin the procedure and have limited accuracy and influence on subsequent aspects of the procedure or on the overall outcome of the surgery.\nFurther complicating matters, ocular tissues are not conducive to pre-surgery reference markings or measurements. This is particularly true because most ocular tissues have wet surfaces and internal physical structures that cannot be accessed for direct measurement or marking prior to surgery.\nAdditionally, pre-surgical rinsing, sterilization, or drug administration to the ocular tissues may dissolve, alter or even remove reference markings prior to the initiation of surgery. Similarly, subsequent wiping and contact with fluids, including the patient's body fluids, during the surgical procedure may remove or distort any remaining reference markings. As a result, even the most accurate surgical reference markings may lose any practical effectiveness beyond the initial stages of the surgical procedure.\nAs such, there is a continuing need for effective surgical reference indicia properly aligned with a particular axis, particularly when proper alignment of pre-surgical data is pivotal to satisfactory patient outcome. For instance, proper alignment of pre-surgical data, particularly the true vertical axis of an eye, with the ocular surgery is highly advantageous with intraocular lens (IOL) implantation and orientation of the IOL within the posterior chamber of the eye.\nThe implantation of an IOL is a highly sophisticated surgical procedure that, in the past, has been performed based on partially accurate or even inaccurate visual measurements and estimated alignment with the vertical or horizontal axis of a patient's eye. Past procedures have commonly relied on the measurement of an eye's vertical or horizontal axis prior to surgery and the subsequent inaccurate translation of that measurement to the IOL implantation procedure where the positioning of the measured axis of the eye may have changed. As a result, it is not uncommon for the IOL placement to be improperly aligned with the true vertical axis of the eye, resulting in such side effects as poor visual acuity, double vision, and halos under low ambient light conditions.\nAccordingly, in spite of the ongoing development and the growing sophistication of contemporary ocular surgery, there is a continuing need for the provision of effective surgical reference indicia including information about or which track a patient's natural vertical axis or other important axis of orientation."}
{"text": "FIG. 1a illustrates a conventional disk drive 10 used for data storage. FIGS. 1a, 1b and 1c are not drawn to scale and only certain structures are depicted for clarity. Disk media 50 is attached to spindle motor and hub 20. The spindle motor and hub 20 rotate the media 50 in a direction shown by arrow 55. Head Stack assembly (HSA) 60 includes a magnetic recording head 30 on actuator arm 70 and positions actuator arm 70 by positioning the voice coil motor (VCM) 25 over a desired data track, shown as recording track 40 in this example, to write data onto the media 50.\nFIG. 1b illustrates an enlarged view of area 30 of FIG. 1a. A magnetic recording transducer 90 is fabricated on slider 80. Slider 80 may be attached to suspension 75 and suspension 75 may be attached to actuator arm 70 as shown in FIG. 1c. \nReferring again to FIG. 1b, Slider 80 is illustrated above recording track 40. Media 50 and track 40 are moving under slider 80 in an in-track direction shown by arrow 42. The cross-track direction is shown by arrow 41.\nThe magnetic recording transducer 90 has a leading edge 91 and a trailing edge 92. In this embodiment, the trailing edge 92 of recording transducer 90 is the final portion of magnetic transducer 90 that writes onto the recording track 40 as the media moves under the slider 80 in direction 42.\nFIG. 1c illustrates a side view of the disk drive 10 shown in FIG. 1a. At least one disk media 50 is mounted onto spindle motor and hub 20. HSA 60 comprises at least one actuator arm 70 that carries suspension 75 and slider 80. Slider 80 has an air bearing surface (ABS) facing media 50. When the media is rotating and actuator arm 70 is positioned over the media 50, slider 80 floats above media 50 by aerodynamic pressure created between the slider ABS and the surface of media 50 facing the ABS of slider 80."}
{"text": "This application is a continuation of U.S. patent application Ser. No. 11/233,322, filed Sep. 22, 2005, which is a continuation of U.S. Pat. No. 6,996,198 issued Feb. 7, 2006, and it claims priority under provisional application No. 60/243,811 filed on Oct. 27, 2000.\nThis invention relates to filters and, more particularly, to filter banks for audio applications.\nIn many areas of audio processing, it is desirable to analyze an audio signal in approximately the same time-frequency form as the human ear (i.e., with bandwidths on the order of one Bark) and with a time resolution that is commensurate with the bandwidth of the filter. In addition, it is desirable to process the signals in each of the bands and then reconstruct them in a manner such that when the bands are unmodified, the filter bank has a nearly perfect reconstruction characteristic. Because the signals might be modified, and different bands might be routed to different devices, not only must the filters provide approximately exact reconstruction, they must also prevent aliasing due to the unequal processing, or modification, of adjacent frequency bands. Hence, an oversampled filter bank is required where aliasing introduced due to unequal processing of bands is below the level of human hearing.\nOne application for this kind of filter bank is found for the problem of separating parts of an individual audio signal into its directs and indirect parts for the purpose of rerouting, in real time, the direct and indirect signals to drivers that reproduce them appropriately. In such an application, a filter bank that approximates the critical bandwidths is essential to identifying the part of a signal with direct cues, and the ability of the reconstruction filter bank to prevent substantial aliasing when adjacent bands are added incoherently (as opposed to coherently) is also an absolute requirement. Hence the need for an oversampled critical band filter bank. In applications that require nonuniform division of signal spectrum, iterative cascaded of uniform filter banks are often used. Iterated filter banks, however, impose considerable structure on the equivalent filters, which results in inferior time-frequency localization compared to direct designs. A study of critically sampled nonuniform filter banks has been reported by J. Princen in “The Design of Nonuniform Filter Banks,” IEEE Transactions on Signal Processing, Vol. 43, No. 11, pp. 2550-2560, November 1995. Nonuniform filter banks studied by Princen are obtained by joining pseudo QMF filter bank sections that are nearly perfect reconstruction filter banks based on cosine modulation and the principle of adjacent channel aliasing cancellation. R. Bernardini et al published “Arbitrary Tilings of the Time-Frequency Plane using Local Bases,” IEEE Transactions on Signal Processing, Vol. 47, No. 8, pp. 2293-2304, August 1999, wherein they describe a cosine-modulation-based structure that allows for time-adaptive nonuniform tiling of the time-frequency plane. Despite their many fine features that are relevant to coding purposes, however, these approaches do not have good aliasing attenuation properties in each of the subbands independently. This makes them unsuitable for tasks where processing effects need to be contained within the bands directly affected. Perfect, or nearly perfect, reconstruction properties of these filter banks in the presence of upsampling are also not clear. The pseudo QMF bank, for instance, loses its aliasing cancellation property if the subband components are not critically downsampled.\nOversampled uniform filter banks based on cosine modulation were studied by Bolceskei et al, and reported in “Oversampled Cosine Modulated Filter Banks with Perfect Reconstruction, IEEE Transaction on Circuits and Systems II, Vol. 45, No. 8, pp. 1057-1071, August 1998, but the cosine modulation places stringent aliasing attenuation requirements."}
{"text": "The Global Positioning System (GPS) is a worldwide radio-navigation system formed from a constellation of 24 satellites and corresponding ground stations. Each satellite continually broadcasts its location in space along with the current time from an internal clock. GPS receivers can determine their position to within a few meters by receiving and analyzing signals transmitted from the satellites. To determine its location, a GPS receiver scans for satellite signals until it has acquired signals from three or more satellites. Two-dimensional locations can be determined by analyzing signals from three satellites, and three-dimensional locations can be determined by analyzing signals from four or more satellites. A GPS receiver determines its location by determining its distance from the GPS satellites based on the received signals and then triangulating these distance measurements. Commercial GPS receivers can determine their locations to an accuracy of 10 meters or less with 95 percent reliability.\nAlthough the current GPS has been successful, it has several shortcomings. For example, GPS satellite signals are subject to errors caused by ionospheric disturbances and satellite orbit discrepancies. Ionospheric and tropospheric refraction can slow satellite signals and cause carriers and codes to diverge. Because ionospheric disturbances vary greatly from location to location, these errors are difficult to correct with civilian-type GPS receivers.\nTo account for these errors, the FAA developed a Wide Area Augmentation System (WAAS) to improve the accuracy, availability and integrity of the GPS. The WAAS is based on a network of approximately 25 wide area ground reference stations (WRSs) that are linked to cover a service area including the entire U.S. and some areas of Canada and Mexico. Each of the WRSs has been precisely surveyed so that its exact location is known. Signals from GPS satellites are received and analyzed by the WRSs to determine errors in the signals, including errors caused by the ionospheric disturbances described above. Each WRS in the network relays its data to a wide area master station (WMS) where correction information is computed. The WMS calculates correction messages for each GPS satellite based on correction algorithms and assesses the overall integrity of the system. The correction messages are then uplinked to a pair of Geostationary Communication Satellites (GEOs) via a ground uplink system (GUS). The GEOs then broadcast the messages on the same frequency as GPS (L1, 1575.42 MHz) to GPS receivers within the coverage area of the WAAS. GPS receivers may then utilize the WAAS correction data to correct for GPS satellite signal errors caused by ionospheric disturbances and other inaccuracies. The communications satellites also act as additional navigation satellites for the GPS receivers, thus, providing additional navigation signals for position determination.\nOne type of information that is included in the correction messages from the GEOs is ionospheric correction data. Ionospheric corrections are broadcast for selected ionospheric grid points generally spaced at 5 degree intervals in both latitude and longitude directions. One approach is to store the correction points in a two dimensional array containing a total of 2,592 elements [(360 degrees longitude divided by 5 degrees) times (180 degrees latitude divided by 5 degrees)]. Many GPS receivers, including, for example, GPS receivers used in avionics applications and portable GPS receivers used for recreational and sport applications have limited memory and processing power and therefore cannot quickly and efficiently store and process all 2,592 ionospheric grid point correction elements. Moreover, even for GPS receivers that have sufficient memory to store all the ionospheric grid point correction data, much of the memory required to do so is wasted because only the correction points within a certain distance of a GPS receiver's current location are relevant to the ionospheric conditions at that location.\nAccordingly, there exists a need for an improved method and system for receiving ionospheric grid point corrections from WAAS satellites. Moreover, there is a need for such a method and system that benefits from the WAAS data while utilizing a minimal amount of memory and system resources."}
{"text": "1. Field of the Invention\nThe present invention relates to permanent magnet DC motors and, more particularly, to electric motors which provide new and improved operation with simultaneous battery regenerating without interfering with the operation of the motor driving system.\n2 . Description of the Related Art\nExemplary type permanent magnet and direct current motor systems are disclosed in U.S. Pat. Nos. 5,117,141; 4,371,801; 4,358,693; 4,237,410; 4,169,983; 3,996,504; 3,633,084; 3,418,506; and 3,334,254.\nElectrical battery systems known in the art suffer from various drawbacks which prevent such systems from achieving maximum performance and reliability. Such problems include a lack of simultaneously charging a long lasting battery energy supply during operation of an electrical battery system. It would, therefore, represent an advance in the art to provide a new improved simultaneous direct current charging system which acts to charge the batteries of a vehicle such as an electric car without having to stop the vehicle."}
{"text": "1. Technical Field\nThe present invention generally relates to flash memory devices, and more particularly, to a method of forming selection transistors as memory transistors in a NAND type flash EEPROM.\n2. Discussion of the Related Art\nNon-volatile memory devices, such as flash memory devices, may be provided in a NOR-type configuration or a NAND-type configuration and can be electrically rewritten and formed with high integration density. NAND-type nonvolatile semiconductor memory devices include a plurality of NAND cell units. Each NAND cell unit is configured by serially connecting a plurality of memory transistors in a column direction between a source and a drain. Selection gate (SG) transistors are connected to at each end of the series-connected memory transistor circuit.\nTwo types of non-volatile memory transistors are floating gate type memory transistors and floating trap (charge trap) type memory transistors. A floating gate type memory transistor includes a control gate and a conductive floating gate that is isolated, by an insulating layer, from a field effect transistor (FET) channel formed in a substrate. Floating gate type memory transistors may be programmed by storing charges as free carriers on the conductive floating gate.\nA floating gate type memory transistor is similar to a standard MOSFET transistor, except that it has two gates instead of just one. One gate is the control gate (CG) like in other MOSFET transistors, but the second gate is a floating gate (FG) that is insulated all around by an oxide insulator. Because the FG is isolated by its insulating oxide layer, any electrons placed on it get trapped there and thus store the information.\nWhen electrons are trapped on the FG, they modify (partially cancel out) an electric field coming from the CG, which modifies the threshold voltage (Vt) of the cell. Thus, when the cell is “read” by placing a specific voltage on the control gate (CG), electrical current will either flow or not flow between the cell's source and drain connections, depending on the threshold voltage (Vt) of the cell. This presence or absence of current is sensed and translated into 1's and 0's, reproducing the stored data.\nFloating trap (charge trap) type memory transistors may include a non-conductive charge storage layer between a gate electrode and a field effect transistor (FET) channel formed in a substrate. Floating trap type memory transistors may be programmed by storing charges in traps in the non-conductive charge storage layer.\nWhen a positive voltage is applied on the gate electrode, electrons are tunneled via the tunneling insulating layer to become trapped in the charge storage layer. As the electrons are accumulated in the charge storage layer, a threshold voltage of the memory transistor is increased, and the memory transistor becomes programmed. In contrast, when a negative voltage is applied to the gate electrode, trapped electrons are discharged to the semiconductor substrate via the tunneling insulating layer. Concurrently, holes become trapped by the tunneling insulating layer. Consequently, the threshold voltage of the memory transistor is decreased, and the memory transistor becomes erased.\nConventional NAND Flash memory strings typically are isolated from other strings by shallow trench isolation (STI), that prevents electrical current leakage between adjacent semiconductor device components, and have three types of transistors which are: the memory transistor transistors (implementing nonvolatile data-storage memory transistors); string select transistors SST; and ground select transistors GST. Typically, in a NAND flash memory device, string selection and ground selection transistors (SSL and GSL) are disposed at the ends of a NAND string and are used to select the NAND string during program, erase and read operations.\nA group of NAND cell units (NAND strings) arranged in a row direction is called a NAND cell block (memory block, MB). The gates of selection transistors SST and GST arranged on the same row are commonly connected to a corresponding one of selection gate lines and the control gates of memory transistors arranged on the same row are commonly connected to a corresponding one of control gate lines. If n memory transistors are serially connected in the NAND cell unit, the number of control gate lines of memory transistors contained in one NAND cell block is n.\nWhen programming data, first, all data items stored in all memory transistors of memory storage cells in the entire memory block (MB) are simultaneously erased. The erase process is performed by setting all of the control gate lines (word lines) of memory transistors in the selected memory block to a low voltage Vss (for example, 0V) and applying high positive voltage Vera (erase voltage, for example, 20V) to a p-type well region in which the memory cell array is formed to discharge electrons in the floating gates into the channel regions. As a result, all data items stored in all memory transistors of memory storage cells in the entire memory block are all set to “1” data. Multiple or all memory blocks can be simultaneously erased.\nAfter the above-described simultaneous data erase step, the data programming process is simultaneously performed for a plurality of memory transistors connected to a selected control gate line. The unit of binary data to be programmed in the memory transistors connected to a selected control gate line is generally defined as one “page” of data. The “page” order in which data is programmed into the memory transistors (pages) in the memory block is based either on a system in which data is programmed in a random order (random programming process) or on a system in which data is sequentially programmed in one direction (sequential programming process). In the sequential programming process, generally, data is programmed in sequential pages in order from the source-side memory transistor.\nIf high positive voltage Vpgm (program voltage, for example, 20V) is applied to a selected control gate line in the simultaneous programming process, electrons are injected from the channel of the memory transistor into the floating gate in the case of “0” data (so-called “0” programming or “0” write). In this case, injection of electrons is inhibited in the case of “1” data (so-called program inhibition, “1” programming or “1” write). Thus, while writing random data into memory transistors of one page, two types of data programming operations are simultaneously performed and it is necessary to control the channel voltage of each memory transistor according to its program data. For example, in the case of “0” data, the channel voltage is kept low so as to apply a strong electric field to the gate insulating film under the floating gate when the program voltage Vpgm is applied to the control gate. In the case of “1” data, the channel voltage is boosted so as to make weak the electric field applied to the gate insulating film and inhibit injection of electrons into the floating gate. If the channel voltage is insufficiently boosted, electrons are injected so that the threshold voltage of the memory transistor to be subjected to the “1” programming process will be changed. This phenomenon is referred as “erroneous programming” or “write error” or “program disturb”. Therefore, in order to realize the programming operation of the NAND type flash EEPROM, it is necessary to suppress variation in the threshold voltage due to erroneous programming within a specified range so as not to cause erroneous operation.\nIn a conventional NAND Flash memory string, the select transistors SST and GST are standard. MOSFET transistors, each having one control gate.\nA leakage current may occur in unselected strings in unselected blocks during a read operation of a selected block, thereby causing read error to occur. Hence, it is necessary to control the leakage current of the select transistor. To this end, conventionally a threshold voltage implant is performed during manufacture in the select transistor region, requiring additional (e.g., mask, implantation) steps in a method of manufacturing the NAND flash memory device.\nFor the purpose of controlling threshold-voltage distributions of programmed memory cells densely and precisely, an incremental step pulse programming (ISPP) mode is often used. According to the ISPP mode, a programming voltage applied to a word line rises stepwise up during repetition of loops of programming cycle. The programming voltage increases by a predetermined step increment (ΔV), also referred to as a “rising rate”. During the programming sequence, a cell threshold voltage of a programmed cell increases at a rate predetermined for each programming loop. Programming of a nonvolatile memory device by means of the ISPP mode is disclosed in U.S. Pat. No. 6,266,270, entitled “Non-Volatile Semiconductor Memory and Programming Method of the Same”. Each programming loop generally is divided into programming and program-verifying periods. In the programming period, memory cells are programmed under a given bias condition as is well known in this art. In the program-verifying period, the memory cells programmed once are verified whether they are conditioned in the target threshold voltages. The programming loops are repeated for a predetermined number of times until all memory cells are completely programmed at the target threshold voltages. As well known, the program-verifying operation is similar to a reading operation, except that read data is not output to external of the device."}
{"text": "(a) Field of the Invention\nThe present invention relates to the art of continuous molten metal casting.\nMore specifically, the invention relates to a device for holding and replacing casting plates of a metallurgical vessel of a casting installation. This plate may be a calibrated plate or a casting tube. These types of plates are usually part of a nozzle comprising a plate connected to a tubular section of varying lengths depending on the applications.\n(b) Description of Related Art\nA device for replacing casting tubes, arranged facing a casting orifice of a metallurgical vessel of a continuous molten metal casting installation is known, particularly from the document EP 0 192 019 A1. Such a device comprises guiding means generally two rails whereon casting tubes can slide, to occupy firstly a standby position, followed by an operating position and finally an evacuation or exit position opposite the standby position. A drive also called actuator or pusher, actuated by a jack or cylinder, is used to push a casting tube from the standby station thereof to the operating station thereof, the tube moved expelling for this reason, to the evacuation station, the worn casting tube situated on the operating station.\nThe casting tube comprises a sliding face wherein a casting channel opens, which is in registry with the metallurgical vessel orifice when the casting tube is at the operating station. The metallurgical vessel orifice generally consists of the casting orifice of the upstream refractory element or the casting orifice of the upstream refractory elements which are in fluid communication. The upstream refractory element is generally rigidly connected to the metallurgical vessel, for example, it is cemented or mortared therein.\nIn the operating position, pushers also called thrusters are arranged extending from the guiding means or rails. These pushers are used to apply a substantially vertical force on two bottom faces of the plate of the casting tube such that the tube sliding face is in tight contact with the face of the upstream refractory element.\nIn some cases, as in WO2004/065041A1 (particularly paragraph 23), the casting tube sliding face is sufficiently large to form, next to the casting orifice, a sealing surface suitable for sealing or closing the metallurgical vessel casting orifice if the casting tube is moved over a distance at least equal to the diameter of the metallurgical vessel orifice. The sealing surface is also called shut-off surface or closure surface.\nThe casting tube present on the operating station can thus adopt two positions:                a casting position, wherein the casting channel thereof is facing the metallurgical vessel casting orifice, and        a sealing position, wherein the sealing surface thereof is facing the metallurgical vessel casting orifice.        \nIn this case, the casting tube may be used not only for casting the molten metal, but also for stopping (interrupting) the casting in the event of an emergency, which is useful for example if another upstream shutdown device is defective.\nThe movement of the casting tube on the guiding means, for instance the rails, thus needs to be controlled selectively, according to whether it is to be moved to the casting position or the sealing position on the operating station, requiring the use of one or more double-stroke jacks. However, such jacks are bulky, heavy and costly. Furthermore, they require the presence of at least two separate hydraulic supplies on the continuous casting floor.\nThe present invention is intended to provide a technical solution for controlling the movement of the casting tube, and more generally of a casting plate, to the casting position or to the sealing position on the operating station, completely automatically, simply and reliably.\nFor this purpose, the present invention relates to a device for holding and replacing a casting plate facing a casting orifice of a metallurgical vessel of a continuous molten metal casting installation, the casting plate being of the type comprising a sliding face wherein a casting channel opens and wherein a sealing surface suitable for sealing the metallurgical vessel casting orifice is formed, said device being of the type comprising a pusher or drive suitable for pushing a casting plate to move it from a standby station to an operating station, a plate on the operating station being suitable for adopting a casting position, wherein the casting channel thereof is facing the metallurgical vessel casting orifice, and a sealing position, wherein the sealing surface is facing the metallurgical vessel casting orifice, the pusher being provided with means for selectively moving same along two strokes:                a short stroke pushing a casting plate to the casting position on the operating station, or        a long stroke pushing a casting plate to the sealing position on the operating station,        \nsaid device being characterised in that it comprises:                a casting plate passage detector between the standby station and the operating station,        a pusher limit switch, controlled by the passage detector and suitable for adopting:        a replacement position corresponding to a casting position, adopted when the detector detects the passage of a casting plate, wherein the limit switch limits the pusher stroke to the short stroke, and        a sealing position, in the other cases, wherein the limit switch allows the pusher to move over the long stroke.        \nBy means of the plate passage detector, if the pusher is actuated and a replacement plate is situated on the standby station, the detector controls the limit switch which is set to the replacement position and limits the stroke of the pusher such that the casting plate is moved to the casting position on the operating station, whereas, if no plate is situated on the standby station, the limit switch allows the pusher to cover the long stroke thereof to push the plate present on the operating station to the sealing position.\nThe operator thus no longer needs to determine whether it is necessary to actuate the pusher for a plate replacement or an emergency stop. The passage detector and the limit switch determine which stroke of the pusher is required automatically.\nIn particular, if the operator actuates the jack without having positioned a replacement plate on the standby station, an emergency stop is required. The device according to the invention thus automatically actuates the jack over the long stroke thereof so that it moves the plate to the sealing position.\nTherefore, the invention provides, besides a simple and economical jack control device, improved safety on the casting site, both for the operator himself, who no longer needs to intervene in the vicinity of the molten metal, and for the entire site, in that the operator can respond more rapidly in the event of an emergency and does not risk to make a mistake.\nThe majority of known devices do not comprise double-stroke jacks, or casting tubes having a sealing surface. When an emergency stop is required, the operator needs to intervene in the vicinity of the molten metal, remove the tube in standby position, replace it by a blank plate and then actuate the jack to move the blank plate into the casting position. The devices equipped with a double-stroke jack and a casting plate comprising a sealing surface are already an improvement as the blank plate and its handling is not necessary anymore. However, they present the drawbacks mentioned above. Double-stroke jacks are bulky, heavy, costly and require the presence of at least two separate hydraulic supplies. These and other problems and disadvantages associated with the prior art are overcome by the invention disclosed herein by providing a simple, economical and safe method for actuating a device. The operator can actuate the jack remotely and very rapidly and seal the casting channel.\nIn one advantageous embodiment, the limit switch is arranged so as to retain the replacement position after detecting the passage of a casting plate, as long as the pusher has not moved back after covering the entire short stroke thereof.\nFor this purpose, a stabilizer is incorporated in the limit switch, so that said limit switch retains the position set by the passage detector, even after the passage detector has stopped detecting the presence of a plate.\nIn one particular embodiment of the invention, the passage detector is a lever which is actuated by the casting plate when it moves from the standby station to the operating station.\nThis means offers the advantage of being simple to produce and reliable in the operation thereof.\nAdvantageously, the limit switch comprises a movable abutment, the pusher or drive comprising a bearing surface suitable for resting on the abutment only when the limit switch is in the sealing position.\nThe lever and the movable abutment may be connected by a ball type connection, which converts the rotation of the lever into a translation movement of the movable abutment. Any other link suitable for transmitting the movement of the lever to the movable abutment could obviously be suitable.\nIn one particular embodiment of the invention, the pusher comprises a rod and the pusher bearing surface is formed by a recess provided in the rod.\nThis embodiment is advantageous in that it is simple to produce and reliable in the operation thereof.\nIn certain embodiments, the recess provided in the rod comprises, opposite the bearing surface, a bevel replacing the movable abutment in the replacement position when the rod moves back after the pusher has covered the entire short stroke thereof.\nAccording to one advantageous embodiment, the device comprises an evacuation or exit station, whereto a worn plate is sent when pushed by a plate pushed to the operating position by the pusher.\nThe invention also relates to an assembly of a casting plate and a device for holding and replacing casting plates wherein the casting plate comprises at least one protrusion or protuberance for interacting with a plate passage detector of a device as described above.\nCladding refractory elements, casting plates or casting tubes by an element such as a metallic casing is known in the art. These casings are well known to those skilled in the art along with the types of materials used to produce said casings. The refractory is preferably contained or cemented into the metallic casing.\nThe invention also relates to a metallic casing for a casting plate of a continuous molten metal casting installation comprising at least one protrusion or protuberance for interacting with a plate passage detector of a device as described above.\nThe casings are generally metallic, particularly made of steel or cast iron. Obviously, any other materials capable of fulfilling the same function could be used. The same applies for the protrusion.\nIn one embodiment, the casing comprises:                a main surface comprising an opening and side edges extending to said main surface and defining the perimeter thereof;        two bearing surfaces substantially longitudinal and intended to slide along the guiding means of the device,        projecting from the main surface, a protrusion extending in the plate sliding direction, the sliding direction being substantially parallel with the longitudinal bearing surfaces.        \nIn one particular embodiment, the casing comprises:                two longitudinal bearing surfaces intended to slide along the rails of the device for guiding the plate,        longitudinal bottom edges parallel with said longitudinal bearing surfaces and,        projecting from at least one of said longitudinal bearing surfaces, a protrusion extending in the plate sliding direction, that is parallel with the longitudinal bearing surfaces.        \nThe bearing surfaces may have various shapes, for example be planar, inclined, or convex. It is simply necessary for them to serve as a support for the casting plate and enable the movement thereof from a standby station to an operating station.\nIn general, the bearing surface is parallel with the plate sliding or replacement direction. In this case, the term “parallel” should be understood in the broad sense, i.e. the bearing surface comprises at least one line segment or generating line parallel with the plate replacement direction. Similarly, an edge or a protrusion is parallel with the bearing surfaces if the edge or protrusion comprises a line segment parallel with the plate replacement direction.\nIn certain embodiments, the casing further comprises one or any combination of any of the following features:                the casing comprises two pairs of opposed side edges as follows: two longitudinal edges and two transverse edges        the two segments respectively parallel to the transverse edges and the longitudinal edges of the casing and comprising the centre of the opening divide the casing into four quadrants; two quadrants being larger        the casing comprises a tubular portion matching and extending from the opening of the main surface        the casing has an overall rectangular outline.        the casing comprises longitudinal bottom edges parallel with said longitudinal bearing surfaces and, projecting from at least one of said longitudinal bearing surfaces, a protrusion extending in the plate sliding direction, the sliding direction being parallel with the longitudinal bearing surfaces        the bearing surfaces are planar        the bearing surfaces are not comprised in the same plane        the casing comprises a pair of opposed side edges, one of which has a first thickness and the second of which has a second thickness greater than said first thickness        the casing is made of cast iron        \nThe protrusion of the casing may be situated on only one side of the metallic casing.\nIn certain embodiments the casing comprises two protrusions wherein each protrusion is situated on either side of the metallic casing, symmetrically in relation to the longitudinal axis of said casing. This configuration is particularly interesting. As explained above, the means for selecting the stroke is located on the pusher. Depending on the casting installation and the space available in the vicinity of the metallurgical vessel, the pusher can be connected on the left side or the right side of the device. In case the vessel comprises a plurality of casting lines each equipped with a device, some may have the pusher on the left side and some others on the right side. Having two protrusions located symmetrically on each side of the plate permit the use of the plate equivalently on all the casting lines, this way, ensuring in all cases the interaction with the passage detector and the correct selection of the stroke.\nIn certain embodiments, the metallic casing protrusion is tapered in the plate sliding direction.\nAdvantageously, the protrusion or each protrusion comprises one or any combination of any of the following features:                the protrusion is formed by a ramp comprising an inclined portion, the inclination being in the plate sliding direction        the protrusion comprises a portion parallel with the bearing surfaces or the bottom longitudinal edges        the protrusion is located outside of the bearing surface        the protrusion is located adjacent the bearing surface        the protrusion is located on the longitudinal sides of a rectangle or outside of a rectangle, the rectangle being formed by the transverse side edges of the casing and the two tangents to the tubular opening parallel to the longitudinal side edges of the casing        the protrusion is located in the two larger quadrants        \nIn view of the high mechanical stress generated on the casing during use as well as the risk of damage of the protrusion(s) or ramp(s) during transport or handling, the casing is in certain embodiments relatively thick and obtained by moulding, e.g. by casting into a mould.\nThe invention also relates to a casting plate for a continuous molten metal casting installation, of the type comprising a sliding face wherein a casting channel opens and wherein a sealing surface capable for sealing a metallurgical vessel casting channel is formed, consisting of:                a refractory defining the casting channel and forming the sliding face,        a metallic casing encasing the refractory in the vicinity of the sliding face,        \ncharacterised in that the metallic casing comprises a protrusion for interacting with a plate passage detector of a device as described above. In certain embodiments, the casting plate comprises a metallic casing as described above.\nAdvantageously, the casting plate comprises one or any combination of any of the following features:                the protrusion of the casting plate projects in the direction opposite the casting plate sliding surface.        the protrusion or each protrusion (30) of the casting plate is formed by a ramp comprised in a plane orthogonal to the sliding face and comprising an inclined portion (30a) and optionally a portion (30b) substantially parallel with the sliding face (19a, 20a).        the plate comprises a refractory tubular extension opposite the sliding face, to extend from the casting channel. The tubular extension may be sufficient to immerse the lower portion thereof in the molten metal mould.        \nThe invention also relates to a method for producing a plate according to the invention comprising the step of assembling a metallic casing and a refractory element. Assembly is performed using known means. In certain embodiments, the refractory is cemented into the metallic casing or assembled by casting a refractory concrete between the refractory element and the casing (cast around). It can also be considered to recover the metallic casing after use and assemble it with a new refractory element."}
{"text": "Time division multiplexing (TDM) is a useful technique for aggregating lower speed data streams into higher speed streams for transmission across a network. Typically, TDM traffic is used in connection-based, fixed bandwidth networks like the SONET data network. In such a network, alignment of data streams is crucial at switching nodes, since data is switched based on its arrival time. TDM switches typically operate based on a static setup that enables desired switching patterns. A new connection is serviced by creating a new data path through the network using control signaling outside the data path—the data path must be fully configured from end to end before any data can be transmitted. By contrast, packet-based or connectionless networks transmit data packets which include the address information required to enable each packet to reach its destination. Connectionless networks can dynamically adopt switching configurations to suit currently available packets. However, connectionless networks are subject to congestion problems if many packets which must share a network path are suddenly simultaneously received. Connection-oriented networks manage congestion perfectly, at the expense of wasted bandwidth if there is no traffic for an established connection.\nTDM networks typically carry multiplexed traffic by allocating a portion of the total available connection time to each connection, based on some regular timing scheme. For example, a single SONET OC-12 connection is a multiplexed aggregate of 12 STS-1 connections. Each STS-1 connection is regularly allocated 1/12th of the total available connection time to serially transmit its next 8 bits (i.e. byte) via an optical connection, as specified in the SONET standard. Each STS-1 connection byte is referred to as a “grain” of that connection. 12 STS-1 connection grains aggregated to make up an OC-12 connection are called a “grain group.” Each grain group is numerically ordered, from the first to the twelfth grain of the corresponding STS-1s. A TDM switch typically accepts a number of multiplexed input connections and switches them, not as the aggregates it receives, but at some lower level, delivering newly reconstituted aggregates at its outputs.\nTDM switches can be characterized by their capability to service arbitrary loads of connection requests (calls). For example, a switch is “non-blocking” if any load of calls can always be mapped successfully from the switch's input ports to the switch's output ports, regardless of any existing call load being serviced by the switch. A switch is “rearrangeably non-blocking” if servicing a new call may require rearrangement (without interruption) of existing calls to enable the switch to service a call that would otherwise be blocked. A “strictly non-blocking switch” can service any new load of calls, provided all of the switch inputs and outputs required to service the load are available (i.e. unused). A switch is “blocking” if there are possible existing call load states for which no input-to-output port mapping rearrangement will enable the switch to successfully map a new call load from the switch's input ports to the switch's output ports.\nStrictly non-blocking switches require hardware proportional to the square of the number of inputs and the number of grains per input, or O(n2g2) where n is the number of inputs and g is the number of grains per input. Hardware complexity can be reduced to O(n2+g2) by using a so-called Clos architecture or its equivalent (Time:Space:Time or Space:Time:Space architectures, for example). But, such hardware savings compromise switching capability: the reduced hardware complexity switch is only rearrangeably non-blocking for unicast calls (i.e. a call in which one ingress grain is scheduled for delivery to one egress port); it is blocking for arbitrary loads of multicast calls (i.e. a call in which one ingress grain is scheduled for delivery to more than one egress port). Moreover, a scheduling algorithm must be used to determine how to route each new call through the reduced hardware complexity switch, taking into account all previously scheduled calls being handled by the switch. Hardware can be added to allow a switch to handle more calls—for example by increasing multicasting capability and eliminating the need to rearrange existing calls. But, regardless of whether a switch is strictly non-blocking for arbitrary loads of multicast calls, rearrangeably non-blocking for restricted widths of multicast calls (i.e. a call in which one ingress grain is scheduled for delivery to less than some predefined maximum number of egress ports), or simply blocking under certain conditions, physical limitations often prevent construction of switches of needed capacities.\nIn order to build larger switches, switching fabrics are formed by connecting switches together in regular patterns. One useful wiring pattern is a 3-stage fabric with 3 columns of switches regularly connected together, as shown in FIG. 1. The first stage (column) of the FIG. 1 fabric has n1 devices, each labeled <a,b> where a=1 and b=1, . . . , n1. The second (center) stage of the FIG. 1 fabric has n2 devices, each labeled <a,b> where a=2 and b=1, . . . , n2. The third stage of the FIG. 1 fabric has n3 devices, each labeled <a,b> where a=3 and b=1, . . . , n3. An output of each first stage device is connected to deliver a maximum input connection load of ti grains to an input of each center stage device. Each first stage device has a total input capacity it grains, and usually it≦n2·ti. An output of each center stage device is connected to deliver a maximum output connection load of to grains to an input of each third stage device. Each third stage device has a total output capacity ot grains, and usually ot≦n2·to.\nIdeally, the center stage devices are designed so that they are non-blocking for whatever level of multicast calls the fabric is required to handle—all multicasting is handled by the center stage. Center stage devices with reduced multicasting capability can be used if corresponding degradation in multicast capability is acceptable. The first and third stage devices are designed so that they are non-blocking for unicast calls. The FIG. 1 fabric has an input speedup factor is=n2·ti/it, an output speedup factor os=n2·to/ot, a total input capacity of n1·it calls and a total output capacity of n3·ot calls. In some applications, is and os simultaneously have values greater than 1. In other applications, is or os, or both, simultaneously have values less than 1. Typically, a fabric's speedup is the same into and out of the center stage, in which case the fabric's overall speedup factor s=is=os. But, s can be expressed in different ways, one common example being s=min(is,os). Of course, if input capacity exceeds output capacity, then some switching capacity is unused. Conversely, if output capacity exceeds input capacity, additional output capacity can be “created” by multicasting calls to achieve full switching capacity usage. Although FIG. 1 shows only one connection between any two devices, persons skilled in the art will understand that each depicted connection represents aggregate connective capacity between two devices. In such fabrics, it is common to have multiple physical links connecting two switches. ti and to thus represent, respectively, total connection load capacity of all physical links into and out of the center stage.\nThe FIG. 1 fabric's aggregate bandwidth is equal to the lowest bandwidth of any of the 3 columns of switches. In this way, fabrics of arbitrary size can be built from any smaller sized switches. The 3-stage fabric is a conceptual organization that can take many real forms. For example, the first and third stage switches of a conceptual 3-stage fabric are often a single physical device. This allows a fabric's first and third stages to be located at a convenient ingress/egress point (physical ingress/egress connectors are often co-located). 3-stage fabrics are somewhat equivalent to Clos networks, in that they are not strictly non-blocking without impractical amounts of speedup, require a scheduling algorithm to derive a switch setting for scheduling a newly received call, and may require rearrangement of existing (i.e. previously scheduled) calls to facilitate scheduling of a newly received call.\nFor each newly received call, the scheduling algorithm examines the new call and the existing switch setting used to service the switch's existing call load, to determine a new switch setting that will preserve all existing end-to-end calls while permitting the new call to be scheduled. Here, “preserve” means that the inputs and outputs of each previously scheduled call are maintained in the new switch setting, although the internal path through the 3-stage fabric may differ (i.e. be rearranged). The optimal new switch setting is one that rearranges the fewest existing calls and leaves the fabric's latent (i.e. unused) switching capability evenly balanced among all switches. The problem of determining optimal switch settings for a 3-stage fabric is “NP-complete,” meaning that there is no known way to solve the problem other than generating all possible switching schedules and examining them to determine which one is optimal.\nFIG. 2 depicts a simple 3-stage switching fabric for which n1=n2=2 and n3=3. The first stage has n1=2 switches, each of which is capable of receiving 3 grains on each of 4 input ports and outputting 6 grains on each of 2 output ports. The second stage has n2=2 switches, each of which is capable of receiving 6 grains on each of 2 input ports and outputting 4 grains on each of 3 output ports. The third stage has n3=3 switches, each of which is capable of receiving 4 grains on each of 2 input ports and outputting 2 grains on each of 4 output ports.\nMore particularly, each output port of each first stage device is connected to deliver ti=6 grains to an input port of each center stage device; and, each output port of each center stage device is connected to deliver to=4 grains to an input port of each third stage device. Each first stage device has a total input capacity it≦n2·ti=2·6=12 grains. Each third stage device has a total output capacity ot≦n2·to=2·4=8 grains. The FIG. 2 fabric has an input speedup factor is=n2·ti/it=2·6/12=1, an output speedup factor os=n2·to/ot=2·4/8=1, an overall speedup factor s=is·os=1·1=1, a total input capacity n1·n2·ti=2·2·6=24 calls and a total output capacity n3·n2·to3·2·4=24 calls.\nIn the particular situation depicted in FIG. 2, switch <1, 1> has received at its input port 1 an input signal consisting of grain A assigned to a first ingress timeslot, nothing assigned to a second ingress timeslot, and grain B assigned to a third ingress timeslot. At its input port 2, switch <1, 1> has received an input signal consisting of nothing assigned to the first ingress timeslot, grain C assigned to the second ingress timeslot, and grain D assigned to the third ingress timeslot. At its input port 3, switch <1, 1> has received an input signal consisting of grain E assigned to the first ingress timeslot, and nothing assigned to the second or third ingress timeslots. At its input port 4, switch <1, 1> has received an input signal consisting of nothing assigned to the first ingress timeslot, grain F assigned to the second ingress timeslot, and grain G assigned to the third ingress timeslot. Switch <1, 1> rearranges its input signals by reassigning (unicasting) grains A and B to its output port 1's egress timeslots 3 and 5 respectively, reassigning grain C to output port 2's egress timeslot 1, reassigning grains D and E to output port 1's egress timeslots 6 and 2 respectively, and reassigning grains F and G to output port 2's egress timeslots 3 and 4 respectively.\nSwitch <1,2> has received at its input port 1 an input signal consisting of grain H assigned to the first ingress timeslot, grain I assigned to the second ingress timeslot, and grain J assigned to the third ingress timeslot. At its input port 2, switch <1,2> has received an input signal consisting of grain K assigned to the first ingress timeslot, and nothing assigned to the second or third ingress timeslots. At its input port 3, switch <1,2> has received an input signal consisting of grain L assigned to the first ingress timeslot, grain M assigned to the second ingress timeslot and nothing assigned to the third ingress timeslot. At its input port 4, switch <1,2> has received an input signal consisting of nothing assigned to the first or second ingress timeslots, and grain N assigned to the third ingress timeslot. Switch <1,2> rearranges its input signals by reassigning (unicasting) grains H and I to its output port 2's egress timeslots 3 and 6 respectively, reassigning grain J to output port 1's egress timeslot 1, reassigning grain K to output port 2's egress timeslot 5, and reassigning grains L, M and N to output port 1's egress timeslots 5, 4 and 3 respectively.\nOutput port 1 of switch <1, 1> is connected as aforesaid to deliver ti=6 grains to input port 1 of switch <2, 1>, output port 2 of switch <1, 1> is connected to deliver ti=6 grains to input port 1 of switch <2,2>, output port 1 of switch <1,2> is connected to deliver ti=6 grains to input port 2 of switch <2, 1>, and output port 2 of switch <1,2> is connected to deliver ti=6 grains to input port 2 of switch <2,2>.\nSwitch <2,1> rearranges its input signals by 2-casting grain A to its output port 2's egress timeslot 1 and to its output port 3's egress timeslot 4, unicasting grain B to output port 2's egress timeslot 4, unicasting grain D to output port 1's egress timeslot 2, 3-casting grain E to its output port 1's egress timeslot 1 and to its output port 2's egress timeslot 2 and to its output port 3's egress timeslot 3, 2-casting grain J to its output port 2's egress timeslot 3 and to its output port 3's egress timeslot 2, unicasting grain L to output port 3's egress timeslot 1, and unicasting grains M and N to output port 1's egress timeslots 4 and 3 respectively.\nSwitch <2,2> rearranges its input signals by 3-casting grain C to its output port 1's egress timeslot 2 and to output port 2's egress timeslot 3 and to output port 3's egress timeslot 1, 2-casting grain F to output port 1's egress timeslot 1 and to output port 2's egress timeslot 1, unicasting grain G to output port 2's egress timeslot 2, 2-casting grain H to output port 1's egress timeslot 4 and to output port 3's egress timeslot 3, and unicasting grains I and K to output port 3's egress timeslots 4 and 2 respectively.\nOutput port 1 of switch <2, 1> is connected as aforesaid to deliver t0=4 grains to input port 1 of switch <3, 1>, output port 2 of switch <2, 1> is connected to deliver to=4 grains to input port 1 of switch <3,2>, output port 3 of switch <2, 1> is connected to deliver to=4 grains to input port 1 of switch <3,3>, output port 1 of switch <2,2> is connected to deliver to=4 grains to input port 2 of switch <3,1>, output port 2 of switch <2,2> is connected to deliver to=4 grains to input port 2 of switch <3,2>, and output port 3 of switch <2,2> is connected to deliver to=4 grains to input port 2 of switch <3,3>.\nSwitch <3, 1> rearranges its input signals by reassigning (unicasting) grains F and M to its output port 1's egress timeslots 1 and 2 respectively, reassigning grains D and N to output port 2's egress timeslots 1 and 2 respectively, reassigning grain H to output port 3's egress timeslot 1, and reassigning grains C and E to output port 4's egress timeslots 1 and 2 respectively.\nSwitch <3,2> rearranges its input signals by reassigning (unicasting) grains F and A to its output port 1's egress timeslots 1 and 2 respectively, reassigning grains G and C to output port 2's egress timeslots 1 and 2 respectively, reassigning grains E and B to output port 3's egress timeslots 1 and 2 respectively, and reassigning grain J to output port 4's egress timeslot 2.\nSwitch <3,3> rearranges its input signals by reassigning (unicasting) grains H and J to its output port 1's egress timeslots 1 and 2 respectively, reassigning grains I and C to output port 2's egress timeslots 1 and 2 respectively, reassigning grains K and L to output port 3's egress timeslots 1 and 2 respectively, and reassigning grains A and E to output port 4's egress timeslots 1 and 2 respectively.\nTo be entirely unambiguous, a call must be specified in terms of its associated input device, port and timeslot; and, its associated output device(s), port(s) and timeslot(s) parameters. A unicast call can be specified as X:(w→x), a 2-cast call can be specified as X:(w→x, y) and a 3-cast call can be specified as X:(w→x, y, z) where X is the label of the call that originates at <stage, device, port, timeslot>w and is received at <stage, device, port, timeslot> x, y and z respectively. Call labels are analogous to grains. Thus, in FIG. 2, each letter-labeled grain represents a different call. The calls currently scheduled on the FIG. 2 fabric are:\nA:(<1,1,1,1>→<3,2,1,2>, <3,3,4,1>)\nB:(<1,1,1,3>→<3,2,3,2>)\nC:(<1,1,2,2>→<3,1,4,1>, <3,2,2,2>, <3,3,2,2>)\nD:(<1,1,2,3>→<3,1,2,1>)\nE:(<1,1,3,1>→<3,1,4,2>, <3,2,3,1>, <3,3,4,2>)\nF:(<1,1,4,2>→<3,1,1,1>, <3,2,1,1>)\nG:(<1,1,4,3>→<3,2,2,1>)\nH:(<1,2,1,1>→<3,1,3,1>, <3,3,1,1>)\nI:(<1,2,1,2>→<3,3,2,1>)\nJ:(<1,2,1,3>→<3,2,4,2>, <3,3,1,2>)\nK:(<1,2,2,1>→<3,3,3,1>)\nL:(<1,2,3,1>→<3,3,3,2>)\nM:(<1,2,3,2>→<3,1,1,2>)\nN:(<1,2,4,2>→<3,1,2,2>)\nFIG. 3A shows the FIG. 2 fabric carrying calls:\nA:(<1,1,1,1>→<3,1,1,1>1)\nB:(<1,1,1,2>→<3,1,1,2>)\nC: (<1,1,1,3>→<3,1,2,1>)\nD:(<1,1,2,1>→<3,1,2,2>)\nE:(<1,1,2,2>→<3,2,1,1>)\nF:(<1,1,2,3>→<3,2,1,2>)\nG:(<1,2,1,1>→<3,2,2,1>)\nH:(<1,2,1,2>→<3,2,1,2>)\nI:(<1,2,1,3>→<3,2,3,1>)\nJ:(<1,2,2,1>→<3,2,3,2>)\nK:(<1,2,2,2>→<3,3,1,1>)\nL:(<1,2,2,3>→<3,3,1,2>).\nSuppose that a new call M:(<1,1,3,1>→<3,2,4,1>) is to be scheduled on the FIG. 3A fabric. That new call cannot be routed through switch <1, 1> to switch <3,2> even though input port 3 of switch <1, 1> is unused and output port 4 of switch <3,2> is unused. This is because the connection between switch <1, 1> and switch <2, 1> is already carrying its maximum of ti=6 grains and the connection between switch <2,2> and switch <3,2> is already carrying its maximum of to=4 grains. Similarly, a new call N:(<1,2,3,1>→<3,1,4,1>) cannot be routed through switch <1,2> to switch <3, 1> even though input port 3 of switch <1,2> is unused and output port 4 of switch <3,1> is unused, because the connection between switch <1,2> and switch <2,2> is already carrying its maximum of ti=6 grains and the connection between switch <2, 1> and switch <3, 1> is already carrying its maximum of t0=4 grains. Potentially available switching capacity via unused input ports 3, 4 of switch <1, 1> to unused output port 4 of switch <3,2> and similar capacity from unused ports of switch <1,2> to switch <3, 1> is thus wasted, reducing the fabric's overall scheduling capability unless the existing call load can be rearranged to improve the fabric's scheduling capability. Such rearrangement is not always possible. It can thus be seen that even if a 3-stage fabric is constructed of strictly non-blocking switching devices, the overall fabric can be blocking, due to the foregoing connection allocation problem.\nFIG. 3B shows the FIG. 2 fabric carrying the same call load as in FIG. 3A, but with the load evenly balanced between each first and second stage switch and between each second and third stage switch. Specifically, the load between switch <1, 1> and switch <2, 1> is 3 grains, as is the load between switch <1, 1> and switch <2,2>, the load between switch <1,2> and switch <2, 1> and the load between switch <1,2> and switch <2,2>. The load between switch <2, 1> and switch <3, 1> is 2 grains, as is the load between switch <2,2> and switch <3, 1>; the load between switch <2, 1> and switch <3,2> is 3 grains, as is the load between switch <2,2> and switch <3,2>; and, the load between switch <2, 1> and switch <3,3> is 1 grain, as is the load between switch <2,2> and switch <3,3>. Thus, although the FIG. 3B fabric is carrying the same call load as the FIG. 3A fabric, the FIG. 3B fabric can carry the new call M:(<1,1,3,1>→<3,2,4,1>) whereas the FIG. 3A fabric cannot.\nFor example, because the load between switch <1, 1> and switch <2, 1> is less than ti=6 grains and the load between switch <2, 1> and switch <3,2> is less than to=4 grains, call M:(<1,1,3,1>→<3,2,4,1>) can be routed through switch <1, 1> to switch <2, 1> and thence to switch <3,2>. Alternatively, because the load between switch <1, 1> and switch <2,2> is less than ti=6 grains and the load between switch <2,2> and switch <3,2> is less than to=4 grains, call M:(<1,1,3,1>→<3,2,4,1>) could be routed through switch <1, 1> to switch <2,2> and thence to switch <3,2>. Similarly, call N:(<1,2,3,1>→<3,1,4,1>) could be routed through switch <1,2> to switch <2, 1> and thence to switch <3,2>; or, through switch <1,2> to switch <2,2> and thence to switch <3,2>.\nThis invention facilitates call load balancing to improve a switching fabric's connection scheduling capability."}
{"text": "As those of skill in the art can appreciate, an automated pipetting system is generally a device that performs programmed transfers of liquid between preselected groups of containers. Laboratory manuals tools called pipettes, used by lab workers and employed in robotic systems alike, are commonly used in application for molecular biology, analytical chemistry, medical tests, and other fields, to transport a measured volume of liquid. As those of skill in the art can appreciate, while such conventional systems for testing are substantially well developed, they are generally engineered to handle liquid samples and are not particularly well adapted for other materials.\nFor example, there are other applications that require the capture of solid specimens for archiving, profiling, monitoring, testing, or other purposes. A common example is in the food and environmental sciences, where samples of solid material, such as the intact muscle mass of an animal, is desired to be sampled. Another is in a medical setting, where a biopsy might be desired from a medical sample. With solid samples, the sample texture can vary significantly, such as for example having soft, hard, dry, moist, oily, and/or viscous forms. These and other characteristics create substantial difficulties with known collection systems and apparatus.\nIn addition to just the physical issues with collection of samples, there are also additional requirements such as turn-around time per sample, convenience of use, or independent sampling events that are important in the context of the hardware solution. As such, the hardware solid sampling solution should be reasonably affordable for the relevant application, relatively easy to assemble and use by skilled persons, and preferably disposable to allow isolated sampling incidents and potentially very short lag-time to process multiple samples back-to-back. Currently, certain biopsy punches are available for solid samples, but most of the foregoing issues are not addressed by such hardware, including the problem that they are not particularly ideal for quick turnaround time applications, lacking good ways of ensuring the captured sample is retained, and not providing an easy manner of releasing the captured samples.\nAccordingly, it would be desirable to provide systems, methods, and modes for sampling of solid or substantially solid materials for food, environmental and related sciences, and for the biological, chemical, medical and related fields."}
{"text": "This invention relates to a tee ball stand. More specifically, this invention relates to a tee ball stand which is self-righting to a vertical orientation in the event the ball holding structure is inadvertently struck instead of the ball itself.\nTee ball stands are characteristically used in the process of teaching young children to hit a ball with a bat. A ball is placed on the end of a pole near the youngster who can then strike at the stationary ball by swinging a bat instead of the more difficult task of attempting to hit a moving ball. With limited experience or limited coordination, the youngster may occasionally strike the pole holding the ball, rather than the ball itself. This can be expected as part of the learning process. As a result, various solutions have been proposed for the safety of the child and for repositioning the tee ball stand in the event of inadvertent and errant bat swings.\nU.S. Pat. Nos. 6,045,462, 6,551,204 and 7,226,372 all relate to either flexure or tilting of the vertical ball support pole by some yielding mechanism. Owen U.S. Pat. No. 6,238,307 discloses a helical spring as a shock absorbing element which directly connects the base and the ball support pole of the tee ball stand.\nIn practice, however, the latter mentioned solution suffers from a number of drawbacks for young children. With the spring member of Owen directly interposed in the ball support member as a shock absorbing element, it is too strong to effectively yield when struck by a bat wielded by a youngster who will still feel the shock of impact through the bat. If, on the other hand, the spring member of Owen is a weaker spring so the youngster will not feel the shock of impact, then the spring will not serve to be self-righting when struck by an errant bat swing.\nTherefore, a need remains in the field of youth sports for a safe tee ball stand that will easily yield when struck and still be self-righting in such situations. The primary objective of this invention is to meet this need."}
{"text": "The lens of the human eye is located centrally behind the pupil and is protected by the cornea. In the normal eye, the lens is clear and substantially symmetrical, with opposed convex surfaces defining generally spherical sections. The lens and the cornea cooperate to focus light on the retina. The retina in turn cooperates with the nerves and the brain, so that light impinging on the retina is perceived as an image.\nThe light refraction which takes place in the cornea and the lens translates into an optical correction of about 60 diopters, with the cornea accounting for about 40 diopters and the lens accounting for about 20 diopters. Other refracting structures also are present in the eye, but are disregarded to simplify the subject explanation.\nA cataract is a condition where the normally clear lens of the eye becomes progressively opaque. This opacification typically occurs over an extended period of time, and the amount of light which passes through the lens decreases with increasing degrees of opacity. As the ability of the cataract lens to transmit light decreases, the ability of the eye to perceive images also decreases. Blindness ultimately can result. Since there are no known methods for eliminating the opacity of a cataract lens, it generally is necessary to surgically remove the opaque lens to permit the unobstracted passage of light through the pupil to the retina. The cataract lens is removed through a generally horizontal incision made at the superior part of the juncture where the cornea and the sclera meet. It is estimated that about 500,000 cataract lenses will be surgically removed in the United States during 1983.\nOnce the cataractous lens has been surgically removed, light can be readily transmitted through the pupil and toward the retina. However, as noted above, the lens of the eye performs a significant light focusing function. Consequently, with the lens removed, the optical system of the eye is left about 20 diopters \"short\", and light is no longer properly focused on the retina. Eyeglasses, contact lenses and intraocular lenses are the three types of optical aids that commonly may be employed after cataract surgery to refocus the light on the retina.\nEyeglasses include lenses which are spaced from the cornea of the eye. The air space between the lens and the cornea causes an image magnification of more than 7%. Unfortunately, the brain cannot assimilate this magnification, and as a result an object appears double. This is a particular problem if the individual had only one cataract eye. Eyeglasses also substantially limit peripheral vision.\nContact lenses rest directly on the cornea of the eye, thus eliminating the air space. As a result, there is a much smaller image magnification with contact lenses than there is with eyeglasses, and the brain typically can fuse the images perceived by an eye with a contact lens and one without. Contact lenses, however, are less than perfect. For example, contact lenses are quite fragile and can be easily displaced from their proper position on the cornea. Additionally, the lenses must be periodically replaced because of protein build-up on the surface of the lens which can cause conjunctivitis. Furthermore, many of the elderly people who require cataract operations do not have the required hand coordination to properly remove or insert the lens.\nIntraocular lenses first became available as optical aids to replace removed cataract lenses in about 1955. These lenses are placed in the eye, and thus closely simulate the optics of the natural lens which they are replacing. Unlike eyeglasses, there is virtually no image distortion with a properly made and placed intraocular lens. Also, unlike contact lenses there is no protein build up on intraocular lenses and the lenses require no care by the patient. Even though intraocular lenses have come into use relatively recently, it is estimated that they will be used for optical correction in approximately 80% of the 500,000 cataract operations anticipated in the United States during 1983.\nThe prior art intraocular lens typically is either of plano-convex construction or double convex construction, with each curved surface defining a spherical section. The lens is placed in the eye through the same incision which is made to remove the cataract lens. As noted above, this incision typically is made along the superior part of the eye at the juncture of the cornea and the sclera. The recipient of an intraocular lens typically will have clear vision with normal peripheral vision if there is no astigmatism. However, in virtually all instances the surgery itself induces astigmatism. Postoperative induced astigmatism is attributable to the healing characteristics of the eye adjacent the incision through which the cataract lens is removed and the intraocular lens is inserted. More particularly, the sutured incision in the eye tends to heal more slowly and less completely as compared to incisions in the skin. For example, a sutured incision in skin typically heals in five to seven days, whereas a comparable incision in the eye may take eight weeks to a year to properly heal depending on the method of suturing. This slow healing rate typically is attributable to the nature of the eyes' tissue, poor vascularity and topical cortisone use after surgery. During the period when the eye is healing, the sutured area tends to spread and thus the prior spherical cornea is made other than spherical. Since the incision always is generally horizontally aligned, the spreading is always along the vertical meridian. Consequently, the optical system of the eye which previously had been spherical becomes \"toric\" with the vertical meridian of the optical system providing a different optical power than the horizontal meridian. This nonspherical configuration of the optic system is generally referred to as \"astigmatism\".\nIt now is known that the degree of this induced astigmatism varies according to the type of sutures used and the suturing technique and the technical skill and care employed by the surgeon. For example, the use of a fine nylon suturing material typically results in a smaller deviation from sphericity than the use of silk or absorbable sutures. Generally, the induced astigmatism varies from 0.5 to 5 diopters. Additionally, the astigmatism resulting from the operation is always in the vertical meridian. The induced astigmatism typically is corrected by prescription eyeglasses where the vertical meridian has more power than the horizontal meridian.\nIt now has been discovered that although the amount of induced astigmatism is subject to several variables, there is little variation in the astigmatism induced by each surgeon. Thus, a surgeon who continually uses the same type of suturing material in the same manner can predict with reasonable accuracy the amount of astigmatism that he or she will induce.\nIn view of the above, it is an object of the subject invention to provide an intraocular lens for eliminating or reducing astigmatism.\nIt is another object of the subject invention to provide an intraocular lens which eliminates or reduces the astigmatism induced by cataract surgery.\nIt is a further object of the subject invention to provide an intraocular lens of toric configuration.\nIt is an additional object of the subject invention to provide a precisely manufactured toric lens with the variation in meridians reflecting the astigmatism that the surgeon anticipates inducing.\nIt is still another object of the subject invention to provide a toric intraocular lens constructed to be properly and easily aligned when placed in the eye."}
{"text": "Beam Profile Reflectometry (BPR) and Beam Profile Ellipsometry (BPE), described in U.S. Pat. Nos. 4,999,014 and 5,042,951 respectively, are established methods for measuring the thicknesses of thin films and coatings deposited upon flat substrates. Each of these techniques makes use of data contained in the cross-sectional profile of a laser beam which has been reflected from the sample under test. As described in each of these patents, each technique relies on the sample under test being (a) perfectly flat and (b) aligned so that an incident probe beam is focused substantially normal to the surface of the sample.\nIn practice, known implementations of BPR rely on a very high degree of accuracy in achieving this alignment, since only when the alignment is perfect are the beam profiles obtained symmetrical, and this in turn is an assumption made in the subsequent analysis. To this end, the present inventor has devised a sophisticated five-axis stage arrangement (described in WO2008/119982) which, when used in combination with optical feedback from a BPR system, can be used to align a sample (including samples with complex shapes) to an arbitrarily high degree of accuracy.\nHowever, the requirement for this alignment (whether achieved statically for flat samples by very accurate construction of supporting hardware, or dynamically for curved surfaces using a stage such as that referred to above) introduces significant cost and complexity into a BPR or BPE system.\nMore specifically, BPR and BPE tend to rely on the ability to fit a calculated optical interference fringe pattern to an experimentally-measured interference fringe pattern obtained from appropriate optical hardware. Where the sample being measured is misaligned, an essential part of this process is to be able to ascertain the nature of the misalignment (i.e. the actual orientation of the sample), which in practice means locating the centre of the beam profile, which, in the case of BPR and BPE comprises an interference fringe pattern. As disclosed in the Applicant's earlier application WO2008/119982, misalignment of the sample causes the image of the reflected laser beam's profile, after it has been refracted through the main objective lens, to become asymmetrical with more light being reflected in the direction towards which the sample is tilted. In that application, however, no quantitative estimate of the amount of misalignment was made, it being sufficient to know in which direction a motorised axis needed to be moved while an optical feedback loop was provided to monitor the changing alignment state.\nIn addition to the effect of misalignment in the sample's orientation, where the measured data is collected using a coherent laser beam as the light source, the interference fringe pattern is typically contaminated by a large amount of so-called ‘speckle’ noise caused by interference between specularly-reflected light from the sample and scattered light from surface defects and/or dust particles on the system's optical components.\nThe Applicants have found that, despite the existence of many algorithms for filtering noisy data and locating the centre of an interference pattern, none of them have proved suitable for the peculiar characteristics of BPR and BPE fringe patterns where (a) the fringe spacing is not regular but becomes much ‘denser’ as one goes from the centre to edge of the pattern, with the rate-of-change depending on the sample refractive index which is not known beforehand, and (b) misaligned fringe patterns do not consist of perfect concentric circles but of ‘distorted’ circles which are also unevenly illuminated.\nIt is therefore an aim of the present invention to address these issues in order both to smooth out a noisy interference pattern and to accurately locate its centre.\nBy way of further background, FIG. 1 illustrates a typical BPR beam profile in the form of an interference fringe pattern 10 for a misaligned sample. As can be seen, this pattern 10 comprises a set of nominally concentric fringes 12 (however, such fringes are not always present in the interference pattern, for example, where very thin coatings are provided on the sample). It will also be clear from FIG. 1 that the interference pattern 10 comprises numerous ‘asymmetries’. In particular, an essentially random laser speckle pattern is overlaid on the fringe pattern itself, and can cause intensity variations of ˜10% of the fringe intensity, sometimes more. Even for a well-aligned sample, if the light source is plane-polarised then the fringes 12 are not perfectly circular: they have reflection symmetry (in the case of a horizontal or vertical polarisation, the symmetry is top/bottom and left/right) but not rotation symmetry. Also, when there are multiple fringes 12, the fringe 12 spacing decreases rapidly as one goes from the centre of the pattern 10 to the edge, but the rate at which this happens depends upon the very properties of the sample that one is usually trying to measure with the BPR apparatus (for example, the thickness of one or more thin film coatings on the sample). This last feature, in particular, rules out the standard method of filtering an image for speckle noise, namely to produce a two-dimensional Fourier transform and then remove the high-frequency components, because the level of filtering necessary to smooth the widest, innermost fringes 12 is likely to obliterate the outermost fringes 12 altogether. Furthermore, when the sample is misaligned, the fringes 12 become ‘compressed’ in the direction of the misalignment and ‘stretched’ elsewhere, so that even the top/bottom and left/right reflection symmetry is lost and the illumination is such that neither the brightest part of the image, nor its centroid (which can be thought of as the ‘centre of gravity’ of the image where brightness corresponds to mass) correspond to the centre of the fringe pattern 10.\nIt is therefore an aim of the present invention to provide an apparatus and method that addresses at least some of the afore-mentioned problems."}
{"text": "This invention relates to 2-(3',5'-difluoro-4'-cyanophenyl)pyrimidine derivatives, and more particularly to novel liquid crystal compositions including 2-(3',5'-difluoro-4'cyanophenyl)pyrimidine derivatives suitable for use in electro-optical displays.\nLiquid crystal display devices utilize electro-optical effects possessed by liquid crystals. The liquid crystal materials used in these devices have a nematic phase, a cholesteric phase and a smectic phase. The most widely used display system uses liquid crystal materials in a twisted nematic mode.\nLiquid crystal display devices have several advantages. The devices are small in size and can be made thin, and can be driven at low voltage with low power consumption. The liquid crystal material is a light receiving element so that when a liquid crystal display is viewed over a long time, eye strain does not occur.\nIn view of these advantages, liquid crystal display technology has been applied to watches, cameras, electronic counters, audio equipment, automobile dashboard indicators, telephone equipment, measuring devices, and the like. More particularly, liquid crystal display devices are also used presently in personal computers and word processor displays and in other devices including displays which require high resolution and many pixels, including black and white and color pocket televisions, and the like. Thus, liquid crystal display devices continue to attract attention as potentially replacing cathode ray tubes. As a result, liquid crystal display devices are widely used in various areas and it is likely that their use will be broadened further.\nFor practical use, liquid crystal compositions must possess the following characteristics.\n1. The liquid crystal materials must be colorless and thermally, optically, electrically and chemically stable; PA1 2. Have a wide nematic temperature range; PA1 3. A rapid electro-optical response speed; PA1 4. Require a low driving voltage; PA1 5. A steep rise in voltage-light transmittance; PA1 6. The temperature dependency of threshold voltage be small; and PA1 7. A wide visual angle.\nMany liquid crystal materials possess the first of the above-desired properties, however, no single compound satisfies all of the remaining characteristics. Thus, liquid crystal compositions are formed of several different nematic liquid crystal compounds or liquid crystal compositions are obtained by mixing liquid crystal compounds with non-liquid crystal compounds to obtain the desirable properties.\nIn order to decrease the driving voltage of a liquid crystal display device, it is necessary to reduce the threshold voltage. However, the following relationship exist between threshold voltage (V.sub.th), the elasticity constant (K) and dielectric constant anisotropy (.DELTA..epsilon.): EQU V.sub.th .varies.(K/.sub..DELTA..epsilon.)1/2.\nThus, in order to decrease V.sub.th, a liquid crystal compound having large .DELTA..epsilon. and K is required.\nAccordingly, it is desirable to provide an improved liquid crystal material having a low threshold voltage and overcomes disadvantages in present liquid crystal compositions."}
{"text": "The present invention relates to evaporative emission control systems for vehicles and in particular to a purge valve that is adapted to be controlled by the engine management control system for regulating the supply of fuel vapors to the engine intake from the fuel tank vapor recovery system.\nIn order to meet current emission requirements, present day vehicles contain evaporative emission control systems which reduce the quantity of gasoline vapors emanating from the fuel tank of the vehicle. Generally, these systems include a charcoal canister which traps the vapors from the fuel tank, and a purge system which draws the vapors out of the canister and feeds them into the intake system of the engine when the engine is running. The fuel vapors are drawn into the engine intake manifold along with atmospheric air drawn through the canister.\nThe capability of the canister to trap vapors from the fuel tank is greatly dependent upon how thoroughly the vapors are purged from the canister when the vehicle was last operated. Accordingly, it is desirable to purge the canister as much as possible while the engine is running. However, the amount of vapor that can be drawn into the engine at any time is limited by the total airflow into the engine and the accuracy with which the purge flow can be controlled. At high speeds or under high engine loads, high purge flow rates can be easily handled. Under such conditions, however, the manifold vacuum is low which tends to limit the amount of fuel vapors and air which can be drawn from the canister into the engine intake manifold. In addition, when the engine is at idle, the airflow into the engine is low. Therefore, purgining at idle must be precisely controlled to prevent a rough idle. Moreover, due to the varying ratio of air to fuel vapors in the purge system, purging during idle can significantly impact the resulting air/flow ratio of the fuel mixture supplied to the engine. Consequently, purging at idle can easily result in a too rich or too lean fuel mixture causing excessive tailpipe emissions unless purging at idle is limited to low flow rates. Current emissions systems, therefore, do not generally purge the canister at idle to any substantial degree.\nHowever, impending tighter emissions requirements and changes to the EPA testing procedures will require larger capacity canisters and therefore higher capacity purge systems. Moreover, the prospect of on-board refueling vapor recovery systems will only add to these system requirements. Accordingly, it is becoming imperative that such systems not only purge at idle, but that maximum flow rates be increased as well. This, of course, presents conflicting requirements for purge systems. Specifically, in order to purge at idle, the purge flow rate must be fairly low and accurately controlled by the engine control computer which monitors the resulting oxygen content of the exhaust gases from the engine. When a canister is saturated with fuel, and vapor is initially purged, the purge flow is very high in fuel vapor. After most of the fuel vapors are drawn out of the charcoal, the purge flow is almost pure air. Therefore, the purge control valve must be capable of allowing the engine control computer to precesely control small flow rates at idle while correcting the idle fuel-air ratio so that tailpipe emissions are not adversely affected. This type of precise flow control is best accomplished using a relatively small valve.\nOn the other hand, it is desirable to purge at very high flow rates when the engine is operating under high speed or heavy load conditions when it can efficiently consume significant quantities of fuel vapor and air with a minimum effect on fuel air ratios. In order to achieve large flow rates, it is necessary for the purge valve to provide a relatively large flow passage. This requirement, of course, is in direct conflict with the requirement for precise low flow rate control. Specifically, it is believed to be impractical to provide a valve large enough to satisfy the high flow requirements which at the same time is capable of precisely modulating the opening of the valve to meet the low flow requirements.\nAccordingly, it is the primary object of the present invention to provide a two-stage purge control valve that is capable of providing both precise control at low flow rates and high flow capacity at low manifold vacuum pressures. In general, this is accomplished by providing a single assembly having two valves which control separate parallel flow paths. Low flow control is achieved with a small solenoid valve adapted to be driven by a pulse width modulated (PWM) signal from the engine control computer. High flow capacity is provided by a vacuum-controlled valve which opens at low manifold vacuum pressures. Because purge flow comprises a relatively small percentage of total air flow into the engine under the conditions when the high flow stage is open, precise control of the high flow capacity valve by the engine control computer is not required.\nAccordingly, the purge valve according to the present invention allows the full range from 10% to 90% duty cycle control to be used to control low flow rates and opens the high flow valve only when the purge flow comprises a small portion of the total engine intake air flow. Moreover, the high flow valve is adapted to open gradually as engine manifold vacuum pressure decreases, thereby proportioning the purge flow to the total engine intake air flow. In addition, the engine control computer can still adjust the high purge flow rate to a degree by controlling the parallel flow through the PWM solenoid valve.\nIn the preferred embodiment of the present invention, the response and flow capacity of both the low and high flow control valves can be calibrated to meet the requirements of a particular engine family or purge system."}
{"text": "There are certain jigging units with movable sieve (Directory on ore processing. Main processes. Moskow, Nedra, 1983, p.p. 53-54). The sieve is given bow-shaped movement with horizontal displacement towards the place of loading material with the sieve moving down and protruding it forward under moving up. As a result the positive effect is achieved due to simultaneous up going of the whole bed and pushing separating material forward along the sieve, the regulation of the sieve movement is done by replacement of the plates with cranks on the movable disks.\nDisadvantage of the device is insufficient disintegration of the movable bed in the loading part of the sieve and mixing up of disintegrating fractions in the loading part. Thus, the efficiency of material disintegration is low. Due to these reasons as well as some others the type of the devices with the movable sieve did not find wide application.\nThere are jiggers with unmovable sieve, such as “THE WEMCO REMER JIG” (Annotation of the device—“THE WEMKO REMER JIG” WEMCO EGUIPMENT, COAL PLANTS 315C Street, St. Albans, W.Va.).\nThe device has an upper unmovable trough with a sieve and a lower movable one, both are connected by rubber diaphragm along perimeter. The lower trough gets vertical reciprocating movement by special doubled eccentric mechanism. As a result vertical pulsation of under-sieve water is being achieved that is very important for jigging.\nDisadvantages of the machine are complexity of the design due to eccentric unit construction, high inertia of the lower part of the machine that results in high energy-consumption because of the fact that each cycle of jigging needs upraising of the lower trough with the whole amount of water, low efficiency of jigging as there is no simultaneous upraising of the whole bed that is characteristic of all jiggers with unmovable sieve.\nAs a prototype of a jigger, the machine with vibrating of the jigging sieve and water in counter-phase is used (Jigging machine. International Patent Classification B03B 05/16, international publication number: WO 02/13974, international publication date: 21 Feb. 2002 (21 Feb. 2002)). The machine consists of troughs connected hingedly by double-arm levers, the upper having a sieve inside and continuous flexible membrane at the bottom interacting with lower trough. The double-arm levers are hinged upon supports. Reciprocating movement of troughs is achieved by setting a drive in the center of axle of double-arm levers rotation, for example, turning gate hydraulic engine. Thereby reciprocating movement of troughs is taking place both in vertical surface—for jigging performance and in the horizontal one—for transporting processing mineral to the place of unloading out of the jigger. The double-arm levers together with troughs form a system of hinge parallelogram, that results in uniform pulsation of under-sieve water and simultaneous uprising of the whole bed, that's very important for efficiency of jigging.\nBecause of the cinematic scheme and correspondingly the design itself the disadvantages of the machine are that the troughs make reciprocal motion and as result of this the considerable part of the drive capacity is spent to overcome inertia rather than to the useful performance. Besides it result in considerable additional loads over the drive and hinges of the machine that is negatively effects reliability of the machine."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication system, a wireless communication apparatus, and a communication method, and in particular, to a wireless communication system including wireless terminals and an access point or access points.\n2. Description of the Related Art\nAs a wireless LAN, a wireless LAN system (ISO/IEC8802-11:1999(E)ANSI/IEEE Std 802.11, 1999 edition) is known which is based on the IEEE802.11 (an IEEE802.11 system also includes an IEEE802.11a system, an IEEE802.11b system and so on). This wireless LAN system employs, as an encryption method, a method called “WEP (Wired Equivalent Privacy)” and which enables privacy to be ensured as in the case with a wired system. Consequently, the security level of a wireless LAN based on the IEEE802.11 has a WEP mode in which the WEP is applied and a non-WEP mode in which the WEP is not applied.\nPractical wireless LAN products according to the IEEE802.11 can communicate in either the WEP mode in which the encryption method “WEP” is applied or the non-WEP mode in which it is not applied. Further, the WEP mode, in which the WEP is applied, includes a 64-bit encryption mode and a 128-bit encryption mode which have different encryption levels. One of these modes is applied to each of the communication or connection links in the wireless LAN to realize communication. In this case, a higher encryption level means a higher security level and stronger encryption.\nOne form of a wireless LAN according to the IEEE802.11 is a system constructed using a plurality of constitutional units called “basic service sets (BSSs)” each composed of one access point and a plurality of wireless clients connected to this access point.\nStructural elements that connect the BSSs together are called “distribution systems (DSs)”. The access point has a function of connecting to the DS. Information is transmitted between a BSS and a DS via the access point. Accordingly, a terminal can communicate with a terminal belonging to another BSS.\nA terminal belongs to a BSS and requires an authentication and association procedures to be executed between itself and an access point in order to communicate with a terminal belonging to another BSS via the access point. Further, when the terminal attempts to reconnect to another access point, a reassociation procedure is executed.\nFor the wireless LAN specified in IEEE802.11, exchanged frames include control frames used for access control, management frames including a beacon or the like, and data frames for data communication.\nBefore a terminal can transmit or receive a data frame to or from an access point, an authentication and association processes must be executed.\nIn the wireless LAN specified in the IEEE802.11, a terminal inquires of an access point whether or not the WEP as an encryption method is used. That is, the terminal requests the access point to use the WEP. When the access point receives this request and if the WEP is available, authentication frames are transmitted between the access point and the terminal. The WEP can be used on the basis of such transmissions of authentication frames.\nAnother form of the wireless LAN specified in the IEEE802.11 is an independently existing BSS, which is called an “IBSS (Independent Basic Service Set)”. The IBSS corresponds to a communication form in which no access points are provided and in which terminals communicate directly with each other. Further, with the IBSS, neither the association process nor the reassociation process are executed. With the IBSS, data frames can be transmitted without executing any authentication processes between terminals.\nIn this manner, in the conventional wireless LANs, communication data are encrypted in order to ensure security. A connection request sender, e.g. a terminal, requests a connection request receiver, e.g. an access point to use an encryption function (WEP function) for communication. If it is possible to use the WEP function according to this request, the access point, receiving this request, accepts the request and encrypts data communication with the terminal. Further, the connection request sender can also take initiative in determining what security level is used for communication.\nIt is expected that wireless LANs will employ, besides the WEP, a plurality of types of encryption methods with different encryption levels, including those having higher security levels than the WEP. Accordingly, it will be desirable to be able to set detailed security levels according to encryption method types, encryption levels, and the like.\nHowever, in the conventional wireless LANs, the minimum encryption level cannot be set for each BSS in order to ensure security. It is thus impossible to make a system that permits only communication based on encryption with a level equal to or higher than the minimum one. Furthermore, it is disadvantageously impossible to set, for communication, detailed security levels according to encryption method types, encryption strengths, and the like.\nMoreover, the IBSS does not require authentication when a data frame is transmitted. Thus, disadvantageously, non-encrypted data frames may be transmitted within system, thus precluding the security in the system from being ensured.\nFurther, security levels preset for the respective BSSs cannot be individually ensured. Likewise, in DS communication executed among a plurality of BSSs, security levels specified for the respective BSSs cannot be individually ensured."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nMany mobile platforms (such as trains, ships, aircraft and automobiles) employ stowage compartments in a cabin of the mobile platform to enable stowage of passenger items, such as carry-on baggage. With regard to commercial passenger aircraft, increased baggage stowage demands have required the stowage compartments to increase in size and load capacity. In addition, there is a drive to increase passengers “personal space” (i.e., headroom) in the cabin of the aircraft. The desire for increased “personal space” in the cabin has resulted in higher ceilings and the placement of storage compartments higher in the cabins.\nThe increased size and load capacity of the stowage compartments coupled with the higher cabin ceilings and higher stowage compartment placement in the cabins can make it difficult for some passengers to close the door on the overhead stowage compartments. Further, the higher placement of the compartments can make it difficult for a passenger to easily determine the status of the compartment, i.e. fully loaded, locked, or overweight. Accordingly, it would be desirable to have an electronic switch system that is capable of displaying data regarding the compartments, while facilitating easier opening and closing of the compartments."}
{"text": "1. Technical Field\nThe present disclosure generally relates to apparatuses for breaking a glass panel unified with a process table, which break the glass panel formed with scribing lines along the scribing lines by irradiating laser beams, and more specifically, to apparatuses for breaking a glass panel unified with a process table to conduct breaking operations of the glass panel after scribing work of the glass panel as one body with the process table, where the apparatuses are installed on either side of the process table in a cutting head transfer-type laser cutting device.\n2. Background\nIn order to cut a glass panel, a glass panel cutting device having a laser is typically used. A glass panel cutting device using a laser forms scribing lines by irradiating laser beams on both sides of a glass panel to be cut, and employs a method of cutting the glass panel formed with the scribing lines with the use of separate breaking apparatuses.\nOne prior art glass panel cutting device using a laser used a transfer type of process table, in which a process table in which a glass panel is located moves while a laser cutting head is fixed thereto.\nFIG. 1 is a perspective view illustrating an overall configuration of a prior art transfer-type process table laser cutting device.\nReferring to FIG. 1, the prior art transfer-type process table laser cutting device comprises laser cutting heads(10) for generating laser beams and irradiating the laser beams, a process table(20) for receiving a glass panel thereon and moving the panel in an X-axis direction in order to form scribing lines on both sides of the glass panel by moving the glass panel, and breaking apparatuses(30) for cutting the glass panel formed with the scribing lines.\nThe transfer-type process table laser cutting device adjusts a cutting width of the glass panel when the laser cutting heads(10) installed on right and left sides of the laser cutting device move in a Y-axis direction along the cutting width of the glass panel before irradiating the laser beams. After the cutting width of the glass panel is adjusted through movement of the laser cutting heads(10) in the Y-axis direction, the process table(20) where the glass panel is placed moves in the X-axis direction to form the scribing lines in the X-axis direction on both sides of the glass panel.\nAfter the process table(20) moves to form the scribing lines on the glass panel, the process table(20) is located on a front side of the laser cutting device. On a front side of the transfer-type process table laser cutting device, the breaking apparatuses(30) are installed in an upper part so as to cut the glass panel formed with the scribing lines on both sides thereof.\nFIG. 2(a) and FIG. 2(b) show front views illustrating a figure before and after operation of each breaking apparatus used for the prior art transfer-type process table laser cutting device.\nReferring to FIG. 2(a) and FIG. 2(b), each breaking apparatus used for the transfer-type process table laser cutting device comprises a frame(32) for fixing the breaking apparatuses to an upper part of a front side of the laser cutting device, clampers(34) configured to move vertically and disposed on right and left sides of a lower part of the frame(32), and breaking bars(36).\nWhen a process table(20) on which a glass panel(5) formed with scribing lines is located is positioned in the lower parts of the breaking apparatuses(30), the clampers(34) and the breaking bars(36) located under the frame(32) move in a Y-axis direction using a Y-axis linear motor in order to adjust the cutting width of the glass panel(5) along the scribing lines.\nAs a result of the Y-axis movement of the clampers(34) and the breaking bars(36), the clampers(34) are located on the insides of the scribing lines formed on the right and left sides of the glass panel(5), and the breaking bars(36) are positioned on the outsides of the scribing lines formed on the right and left sides of the glass panel(5).\nAfter adjusting the cutting width of the glass panel(5) by the Y-axis movement of the clampers(34) and the breaking bars(36), the clampers(34) and the breaking bars(36) move to the lower part via cylinders. The clampers(34) fix the glass panel(5) to the process table(20) by falling to a surface in contact with an upper side of the glass panel(5), and the breaking bars(36) fall further than the clampers(34) thereby cutting the cutting sections of the glass panel(5) along the scribing lines.\nAccording to the aforementioned transfer-type process table laser cutting device, the breaking apparatuses of the glass panel conduct breaking operations by being installed at separate positions from the laser cutting heads. This can cause an overall layout space occupied by the laser cutting device of the glass panel to be quite large.\nA transfer-type cutting head laser cutting device has been introduced to reduce the layout space required for the transfer-type process table laser cutting device and to increase productivity by reducing the cycle time of a glass panel cutting process. However, the transfer-type cutting head laser cutting device is equipped with a gantry structure for moving laser oscillators and laser irradiation heads in an X-axis direction in an upper part of a process table, which prevent prior breaking apparatuses that apply pressure to a lower part from being installed in upper parts."}
{"text": "This invention relates to steroidal glycosides and methods of using the same, particularly as hypocholesterolemic agents and antiatherosclerosis agents, in mammals.\nMany known products possessing hypocholesterolemic activity are cross-linked synthetic polymer derivatives. For example, cross-linked, water-insoluble, bile-acid-binding polystyrene-based resins, e.g., Cholestyramine.RTM. agents, have a gritty \"mouth-feel\", and thus have poor palatability. In addition, these resin beads typically have a low in vivo efficiency. Thus, the effective hypocholesterolemic dose of these materials is excessive, typically 18-24 grams of formulated product per day. Other known polymers having hypocholesterolemic activity include the natural product chitosan and chitosan derivatives as described in European Application pub. no. 0212145. However, the effective hypocholesterolemic dose of these materials is also high.\nOther known hypercholesterolemia controlling agents include plant extracts such as \"alfalfa saponins\". However, these plant extracts are of variable composition and contain significant amounts of non-useful chemical substances. Due to the variations in composition, it is difficult to set a standard dosage or predict the impurities present. Thus, such extracts are not well suited for use by humans. Furthermore purification of these extracts would be expensive. As an alternative certain synthetically produced, pure \"sapogenin-derived\" compounds e.g., substances compounded from spirostane, spirostene or sterol-derived compounds depress cholesterol absorption more effectively than alfalfa extracts on a weight basis and thus can be administered in reasonable sized doses. Because the chemical compositions of these substances are known and because they can be synthesized at a high degree of purity, they are suitable for use by any warm-blooded animal, including humans.\nHowever, unless administered in massive amounts, pure sapogenins do not significantly inhibit cholesterol's absorption. It is only when compounded with another moiety that sapogenins have the desired effect. Examples of such sapogenin compounds are compounds of tigogenin and diosgenin, particularly glycosides thereof. P. K. Kintia, Iu. K. Vasilenko, G. M. Gorianu, V. A. Bobeiko, I. V. Suetina, N. E. Mashchenko, Kim. Pharm. Zh., 1981, 15(9), 55 discloses 3-O-(.beta.-D-galactopyranosyl)hecogenin and its use as a hypocholesterolemic agent. U.S. Pat. Nos. 4,602,003 and 4,602,005 disclose certain steroidal glycosides, in particular 3-O-(.beta.-D-glucopyranosyl)tigogenin and 3-O-(.beta.-D-cellobiosyl)tigogenin and their use for the control of hypercholesterolemia. 3-O-(.beta.-D-cellobiosyl)tigogenin has superior hypocholesterolemic activity when compared to, for example, cholestyramine. PCT publication WO 93/07167 discloses several steroidal glycosides in particular 3-O-(5-C-hydroxymethyl-L-arabino-hexopyranosyl)-tigogenin and 3-O-(5-C-hydroxymethyl-L-arabino-hexopyranosyl)-diosgenin and their use in the control of hypercholesterolemia.\nRecently commonly assigned PCT publication WO 93/11150 has disclosed a number of steroidal glycosides including 11-ketotigogenyl-beta-O-cellobioside, hecogenin-beta-O-cellobioside, diosgenin-beta-O-cellobioside and their use as antihypercholesterolemic agents. Also commonly assigned PCT publication WO 94/00480, the disclosure of which is hereby incorporated by reference, discloses a variety of steroidal glycosides and their use as antihypercholesterolemic agents.\nAlthough the hypocholesterolemic compounds described above make a significant contribution to the art there is a continuing search in this field of art for improved hypocholesterolemic pharmaceuticals."}
{"text": "Bike racks that clamp the front fork of a bike (a “fork mount”) are desirable particularly for high-end bike frames. These racks avoid the need to clamp the tubes of the frame, and they hold the bike more firmly, substantially avoiding incidental vibration and movement of the bike frame in transit. However, when a fork mount is used to carry a bike, the front wheel that has been removed must be carried somewhere else. Thus, there is sometimes a need for a device for carrying an isolated wheel. Similar needs exist for support vehicles that carry spare wheels for races or organized road tours. This type of device may be referred to as a “wheel carrier, a “front wheel holder,” or a “wheel mount.”\nWheel carriers are inadequate for many purposes. Some wheel carriers are not mountable on crossbars of various shapes and sizes. Some wheel carriers are not versatile enough to carry wheels of different sizes. Some wheel carriers do not hold a wheel sufficiently tightly to minimize vibrations, which may be noisy, annoying, and potentially damaging to the wheel and/or wheel carrier. Some wheel carriers do not adequately protect against theft. Thus, there is a need for an improved wheel carrier."}
{"text": "A hydraulic system may include a variety of hydraulically actuated components, each of which may have different flow and pressure requirements that can vary over time. The hydraulic system may include a pump for supplying a flow of pressurized fluid to the hydraulic components. The pump may have a variable or fixed displacement configuration. Fixed displacement pumps are generally smaller, lighter, and less expensive than variable displacement pumps. Generally speaking, fixed displacement pumps deliver a finite volume of fluid for each cycle of pump operation. The output volume of a fixed displacement pump can be controlled by adjusting the speed of the pump. The pump may be sized to satisfy a maximum flow requirement of the hydraulic system. Closing or otherwise restricting the outlet of a fixed displacement pump will generally cause a corresponding increase in the system pressure. To avoid over pressurizing the hydraulic system, fixed displacement pumps typically utilize a pressure regulator or an unloading valve to control the pressure level within the system during periods in which the pump output exceeds the flow requirements of the hydraulic components. The pressure regulator or unloading valve operates to redirect the excess fluid back to a hydraulic system sump to be re-pressurized by the pump. This method of controlling system pressure and flow may result in a significant reduction in the operating efficiency of the hydraulic system depending on the duration and magnitude of excess pump flow. The hydraulic system may further include various valves for controlling the distribution of the pressurized fluid to various hydraulic components."}
{"text": "This application claims priority to Application No. 10026240.6 which was filed in the German language on May 26, 2000.\nThe invention relates to a method for optical polarization control, and in particular, to endless polarization control of an optical signal.\nFor optical polarization control, polarization transformers are particularly suitable, preferably those of an electro-optical principle of operation. These transformers typically include one or more polarization actuators, the modes of which can be varied endlessly, i.e. without interruption, on a great circle of the Poincarxc3xa9 sphere. Examples of this are known from IEEE J. Quantum Electronics 18(1982)4, pp. 767-711, from IEEE J. Lightwave Techn. 6(1988)7, pp. 1199-1207, from IEEE J. Lightwave Techn. 8 (1990) 3, pp. 438-458, from Electron. Lett. 27(1991)4, pp. 377-379, and from the minutes of the European Conference on Optical Communications 1993, Montreux, Switzerland, pp. 401-404, contribution WeP9.3.\nIn practice, non-ideal behavior of polarization actuators can impede the capability for endless, uninterrupted polarization tracking which is always the aim. In IEEE J. Lightwave Techn. 8(1990)3, pp. 438-458, it was specified how this problem can be solved purely quantitatively, i.e. by adding another electro-optical polarization actuator. More elaborate tests have shown, however, that only one additional polarization actuator is hardly sufficient in practice.\nIn IEEE J. Lightwave Techn. 6(1988)7, pp. 1199-1207, a solution to solving this problem is to use a calibration table in a manner which provides the maximum control speed per control step. This can be disadvantageous, however, in that the constructional length of the polarization transformer normally required is distinctly increased by adding a further polarization actuator to an original electro-optical polarization actuator which sufficient in the ideal case, and is even double in this case. As an alternative, the control voltages could be distinctly increased while keeping the total construction length unchanged. Both methods contradict the usual requirements for low supply voltages and short constructional length for the purpose of avoiding DC drift and optical insertion losses. Using this technique, the further polarization actuator is driven as a function of the parameters of the original one. In the nomenclature used there, a phase angle difference, occurring for d3=0, d4=0, of d2xe2x80x2=atan2(f4(d2), f3(d2)) (atan2=arcfunction, known for example from the programming languages Pascal, C, Matlab) of the further actuator which occurs in the case of input polarization P1, which is circular here, between this and a proportion of output polarization P2 (which is orthogonal thereto and is thus oppositely circular here), is also a function of a corresponding phase angle difference d2 of the original one which is accounted for by a calibration table. Since figure 15 there is a twisted curve, the problem arises that atan2(f4(d2),f3(d2)) and d2 can differ, for example by an odd-number multiple of xcfx80. In these cases, the total delay of the polarization transformer formed by the two polarization actuators is obtained by subtracting the delay d1 of the original one and the delay sqrt (f3(d2){circumflex over ( )}2+f4(d2){circumflex over ( )}2) of the further polarization actuator. That is, d1xe2x88x92sqrt(f3(d2){circumflex over ( )}2+f4(d2){circumflex over ( )}2) for d1 greater than sqrt(f3(d2){circumflex over ( )}2+f4(d2){circumflex over ( )}2). This destructive interplay is the reason why a large constructional length and/or high control voltages are needed in accordance with the prior art. If the method specified is properly performed, atan2(f4(d2),f3(d2)) will change by 4xcfx80 or d2 by 2xcfx80 so that the amount |atan2(f4(d2),f3(d2))xe2x88x92d2 of the difference atan2(f4(d2),f3(d2))xe2x88x92d2 can grow without limits if d2 grows arbitrarily, e.g. by many times 2xcfx80. In particular, the discussed destructive interplay occurs there, for example with d2=0.7*xcfx80, d2=1.3*xcfx80 and d2=1.8*xcfx80.\nSince this phase angle difference (d2 or atan2(f4(d2),f3(d2))) in each case represents an angle coordinate from an eigenmode of the polarization actuator on a great circle on the Poincarxc3xa9 sphere, the words angle coordinate will be used synonymously with this phase angle difference in the text which follows.\nIn IEEE J. Lightwave Techn. 6(1988)7, pp. 1199-1207, an electro-optical polarization transformer for transforming a particular one into any arbitrary polarization state or conversely has been described; such polarization transformers could be abbreviated by GSBA (Generalized Soleil Babinet Analogue). Polarization transformers for transforming any arbitrary polarization state into any arbitrary one will be abbreviated by ER (Elliptical Retarder) in the text which follows. ERs with calibration table(s) for compensating for non-ideal component behavior have not previously been known. This is especially true of ERs that necessitate a short constructional length since, compared with GSBAs, ERs already need twice the constructional length in the ideal case.\nIn Proc. Fourth European Conference on Integrated Optics ECIO 87, Glasgow, Scotland, pp. 115-118, a GSBA is specified with many sections and fixed angle coordinate differences (For example: xcex1*xcfx80/2 between adjacent electrodes, where xcex1=xcex9DEV/xcex9BEAT), which are only dependent on the wavelength but not on the desired polarization transformation, of different sections. Correspondingly, this GSBA also does not have a calibration table for compensating for non-ideal behavior.\nIn one embodiment of the invention, there is a method for endless polarization control of an optical signal. The method includes, for example, converting a polarization of the optical signal using polarization acutuators, with arbitrary, freely and interruption free selectable individual phase angle differences as individual angle coordinates, at least partially into a respective orthogonal, which together form a polarization transformer with a total delay and a total phase angle difference as a total angle coordinate between the polarization and a respective orthogonal, in which an individual angle coordinate is varied as a function of another angle coordinate, wherein the amount (|dxcex61|, |dxcex62|, |dxcex6|, |dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)|) of the variation (dxcex61, dxcex62, dxcex6, dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)) of an individual angle coordinate is limited for sequences of desired polarization transformations.\nIn another aspect of the invention, individual delays support each other in their effects ((\"psgr\"1xe2x80x2 greater than \"psgr\"1, \"psgr\"xe2x80x2 greater than \"psgr\"2), (\"psgr\" greater than \"psgr\"1, \"psgr\" greater than \"psgr\"2)) by the limitation of the amount (|dxcex61|, |dxcex62|, |dxcex6|, |dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)|) of the variation (dxcex61, dxcex62, dxcex6, dxcex6*(2*ixe2x88x92nxe2x88x921)/(nxe2x88x921)).\nIn another aspect of the invention, the another angle coordinate is the actual or attempted total angle coordinate or an attempted individual angle coordinate.\nIn yet another aspect of the invention, a number of individual angle coordinates are varied as a function of the another angle coordinate.\nIn another aspect of the invention, at least in the case of equal attempted individual angle coordinates, one of the variations of individual angle coordinates is equal to the negative of another one of the variations or equal to zero.\nIn another aspect of the invention, at least one of the individual delays is varied as the function of another angle coordinate.\nIn still another aspect of the invention, one of the variations of the individual delays is equal to another one of the variations or equal to zero, at least in the case of equal attempted individual delays.\nIn another aspect of the invention, the method includes applying, to at least two additional polarization actuators, which are located before or after the polarization actuators in the beam path to produce another polarization transformer and form an elliptical retarder.\nIn another aspect of the invention, one of the polarization actuators operates as an electrooptical Soleil Babinet compensator which can convert circular polarizations at least partially into one another.\nIn yet another aspect of the invention, at least two additional polarization actuators producing an additional polarization transformer form an elliptical retarder with the polarization transformer.\nIn another aspect of the invention, at least one of the polarization actuators operates as an electrooptical Soleil Babinet compensator, the polarizations which are at least partially converted being circular polarizations.\nIn another aspect of the invention, at least one of the polarization actuators operates as an electrooptical Soleil Babinet analog, the polarizations which are at least partially converted being TE and TM polarizations."}
{"text": "1. Field\nThe present invention relates to a method for management of a plurality of Communication devices in an access network as described in preamble of claim 1 and an access unit as described in the preamble of claim 10.\n2. Description of Related Arts\nSuch management-method and related device are generally known in the art. The management method, for instance configuration, of communications devices in an Internet Protocol network access network is realised using the Internet Protocol layer for provisioning the layer three connection between a management network element such as configuration server and a to be managed device, required for management purposes.\nSuch a communications network comprises an Access Unit and an Ethernet, ATM or a MPLS aggregation network, wherein the Access Unit, such as a Digital subscriber line access unit, couples the communications device, such as a Digital subscriber line modem, to the aggregation network. The communications device being coupled to the Access Unit over an ATM connection.\nThis communication required for the management of the communications devices is based on the Internet Protocol address further referred to as IP-address, for addressing the management application at the communications device.\nCurrently, the communication required for the management of the communications devices is based on Internet Protocol, where the Internet Protocol addresses of the communication devices are unique addresses.\nIn some situations however, no unique IP-address is assigned to the communications devices (communication devices may have the same IP address) or the communication devices do not have an externally known IP-address. Then it is not possible to use the IP layer for addressing communications devices for management purposes and other purposes.\nThis situation could occur for example due the lack of public IP addresses, or in a network where the uniqueness of the IP addresses assigned to the communications devices cannot be ensured.\nHence, if there is no unique IP-address available for addressing such communications devices communications devices are not reachable at layer three and not able to be managed."}
{"text": "1. Technical Field\nThe present disclosure relates to a home interior monitoring system and a communication control method.\n2. Description of the Related Art\nIn recent years, there has been known a monitoring camera system (for example, refer to Japanese Patent No. 5845453) that includes a monitoring camera, a master device of a fixed phone that is connected to a fixed phone network and can communicate with another fixed phone, and a smartphone that can wirelessly communicate with the master device through a wireless router and wirelessly connects with another mobile phone through a mobile phone network (for example, 3G or 4G).\nIn the monitoring camera system of Japanese Patent No. 5845453, the smartphone, in the case of a wireless router not being registered in the master device, displays a screen prompting a predetermined operation to be performed on the master device and waits for a signal from the master device. The master device directly and wirelessly connects with the smartphone by the predetermined operation and sends a setting request signal to the smartphone. The smartphone, when receiving the setting request signal, wirelessly connects with a wireless router to acquire router setting information from the wireless router and wirelessly connects with the master device to send the router setting information to the master device. The master device registers the wireless router based on the router setting information. Accordingly, when an instruction operation that causes, for example, the smartphone to display an image is performed, the master device receives information as to the instruction operation from the smartphone through the wireless router and can send image data transferred from the monitoring camera to the smartphone through the wireless router.\nIn the monitoring camera system in the related art, data cannot be sent and received between the master device and the smartphone without the wireless router, and convenience is not sufficient. If the master device, instead of the wireless router that can connect to the Internet, operates as an access point (AP) that cannot connect to the Internet, the smartphone is expected to have difficulty in continuing connection with the master device."}
{"text": "1. Field of the Invention\nThe present invention relates generally to attaching bags to objects and more specifically to a magnet system for removable attachment of a bag to a water board, which allows a bag to be removably secured to a water board.\n2. Discussion of the Prior Art\nPatent application no. 2007/0199966 to Korchmar discloses an article storage bag. The Korchmar patent application teaches a handbag with a magnetic closure system for removably retaining a flap against a side of a storage bag body. Patent application no. 2013/0122761 to Macias discloses a stand up paddleboard sweeper. The Macias patent application teaches the attachment of a bag to a paddle board with fasteners.\nAccordingly, there is a clearly felt need in the art for a magnet system for removable attachment of a bag to a water board, which allows a bag to be removably secured to a water board, such as a surf board or paddle board."}
{"text": "1. Field of the Invention\nThe present invention relates to a component processing system and a component processing method for operating an application using a component.\n2. Description of the Related Art\nA component means a software component implemented with business logic and so on for configuring an application. EJB (Enterprise Java Beans (registered trademark)) that is a code proposed by Sun Microsystems, Inc. in US, and so on are examples thereof. The EJB is a standard component architecture for configuring distributed-object-oriented business applications in JAVA (registered trademark) programming language, and it enables configuring distributed applications by combining components that are developed by using tools of different vendors. The use of the EJB allows a developer to be dedicated only to the development of the business logic when developing applications, so that efficient application development is made possible.\nIn the EJB architecture, a container provides an execution environment of the EJB component. Normally, this container is managed by an EJB server. When an EJB client (application) requests the container to create an object of EJB via a standardized interface provided by the container, the container reads description, which is called Deployment Descriptor, defining operational conditions (attribute values of transaction management, security management, and so on) of the component and installs the component based on the operational conditions to create an object. Consequently, an executable component is created (see, for example, Japanese Patent Application Laid-open No. 2002-49496)."}
{"text": "1. Field of the Invention\nThe present invention relates to a pump apparatus and more particularly to a pump apparatus having a valve within a piston for a metering and dispensing device that may be used for food products, the device being very efficient, reliable and simply constructed as well as easy to use and to wash.\n2. Description of the Related Art\nDispensers for food products are known in the food industry. For example, a dispenser for barbecue sauce is constructed with a piston in a cylinder, an intake check valve and an outlet check valve. A major problem with this dispenser is that the food product being transmitted includes chunks which clog the check valves and require frequent cleaning. This procedure is very time consuming and thus the dispenser is not practical."}
{"text": "Wide-bandgap semiconductor devices are based on a semiconductor material with a bandgap of at least 2 eV or at least 3 eV and exhibit lower on-state resistance at high temperatures, lower switching losses and lower leakage currents compared to conventional silicon-based semiconductor devices. Semiconductor devices from wide-bandgap material may include asymmetric transistor cells with stripe-shaped trench gate electrodes that control transistor channels in only one of two opposite longitudinal mesa sidewalls of mesa portions formed from the semiconductor material between neighboring trench gate structures.\nIt is desirable to improve device characteristics of wide-bandgap semiconductor devices including asymmetric transistor cells with trench gates and to further expand the range of applications for such devices."}
{"text": "The present invention relates in general to electrocardiographic (ECG) and oxygen saturation signal recording and more particularly concerns a novel technique for processing the ECG and oxygen saturation signals.\nFor background on ECG and oxygen saturation signal recording reference is made to U.S. Pat. No. 6,125,296.\nA typical prior art approach for measuring oxygen saturation uses either a large nonportable bedside unit or a portable unit with recording capability.\nAn object of this invention is to provide an automatic mechanism for identifying those portions of the recorded pulse oxymetry data signals that are invalid.\nAnother object of the invention is to use pulse oxymetry data signals, which this invention determines as valid, and identify those segments of the data signals that are valid desaturation signals.\nIn one aspect, the invention is a method for processing pulse oxymetry data signals. The method includes recording pulse oxymetry data signals. The pulse oxymetry data signals have a plurality of oxymetry waveforms. The pulse oxymetry data signals correspond to oxygen saturation data signals. The method also includes determining a correlation coefficient between sequential oxymetry waveforms and identifying a valid pulse oxymetry waveform.\nThis aspect of the invention may have one or more of the following features. Determining the correlation coefficient includes storing a first pulse oxymetry waveform segment having a first length in a first buffer, copying the first pulse oxymetry waveform segment from the first buffer to a second buffer, copying a second pulse oxymetry waveform segment having a second length to the first buffer, comparing the first length, the second length and a predetermined value, and determining a correlation length related to the first and second lengths and the predetermined value. Determining a correlation length includes taking the minimum of the first length, the second length, and the predetermined value. Determining a correlation coefficient includes determining a correlation coefficient from the first pulse oxymetry waveform segment and the second pulse oxymetry waveform segment, comparing the correlation coefficient to a threshold value, and filtering out an invalid pulse oxymetry waveform segment that has a correlation coefficient below the threshold value. Filtering out the invalid pulse oxymetry waveform segment includes eliminating pulse oxymetry waveform segments if 75% of the correlation coefficients for the last 6 seconds are above the threshold value of 0.9. The method also includes determining valid oxygen desaturation data signals. Determining valid desaturation signals comprises labeling oxygen saturation signals below a threshold value for a predetermined time as valid desaturation data. The threshold value is 88% and the predetermined time is 300 seconds. Determining valid desaturation signals comprises eliminating artifacts. Determining valid desaturation signals comprises ignoring desaturation signals below the threshold value that do not reach a peak value.\nIn another aspect, the invention is an apparatus for processing pulse oxymetry data signals. The apparatus includes a memory that stores executable instruction data signals and a processor. The processor executes the instruction data signals to record the pulse oxymetry data. The pulse oxymetry data has a plurality of oxymetry waveforms. The pulse oxymetry data signals correspond to oxygen saturation data signals. The processor also executes the instruction data signals to determine a correlation coefficient between sequential oxymetry waveforms and to identify invalid pulse oxymetry waveforms.\nThis aspect of the invention may have one or more of the following features. Determining the correlation coefficient includes storing a first pulse oxymetry waveform segment having a first length in a first buffer, copying the first pulse oxymetry waveform segment from the first buffer to a second buffer, copying a second pulse oxymetry waveform segment having a second length to the first buffer, comparing the first length, the second length and a predetermined value, and determining a correlation length related to the first and second lengths and the predetermined value. Determining a correlation length includes taking the minimum of the first length, the second length, and the predetermined value. Determining a correlation coefficient includes determining a correlation coefficient from the first pulse oxymetry waveform segment and the second pulse oxymetry waveform segment, comparing the correlation coefficient to a threshold value, and filtering out an invalid pulse oxymetry waveform segment that has a correlation coefficient below the threshold value. Filtering out the invalid pulse oxymetry waveform segment includes eliminating pulse oxymetry waveform segments if 75% of the correlation coefficients for the last 6 seconds are above the threshold value of 0.9. The processor also executes the instruction data signals to determine valid oxygen desaturation data signals. Determining valid desaturation signals comprises labeling oxygen saturation signals below a threshold value for a predetermined time as valid desaturation data. The threshold value is 88% and the predetermined time is 300 seconds. Determining valid desaturation signals comprises eliminating artifacts. Determining valid desaturation signals comprises ignoring desaturation signals below the threshold value that do not reach a peak value.\nIn still another aspect, the invention is an article that includes a machine-readable medium that stores executable instruction signals for processing pulse oxymetry data signals. The instruction signals cause a machine to record the pulse oxymetry data signals. The pulse oxymetry data signals have a plurality of oxymetry waveforms. The pulse oxymetry data signals correspond to oxygen saturation data signals. The instructions also cause a machine to determine a correlation coefficient between sequential oxymetry waveforms and to identify invalid pulse oxymetry waveforms.\nThis aspect of the invention may have one or more of the following features. Determining the correlation coefficient includes storing a first pulse oxymetry waveform segment having a first length in a first buffer, copying the first pulse oxymetry waveform segment from the first buffer to a second buffer, copying a second pulse oxymetry waveform segment having a second length to the first buffer, comparing the first length, the second length and a predetermined value, and determining a correlation length related to the first and second lengths and the predetermined value. Determining a correlation length includes taking the minimum of the first length, the second length, and the predetermined value. Determining a correlation coefficient includes determining a correlation coefficient from the first pulse oxymetry waveform segment and the second pulse oxymetry waveform segment, comparing the correlation coefficient to a threshold value, and filtering out an invalid pulse oxymetry waveform segment that has a correlation coefficient below the threshold value. Filtering out the invalid pulse oxymetry waveform segment includes eliminating pulse oxymetry waveform segments if 75% of the correlation coefficients for the last 6 seconds are above the threshold value of 0.9. The instructions also cause a machine to determine valid oxygen desaturation data signals. Determining valid desaturation signals comprises labeling oxygen saturation signals below a threshold value for a predetermined time as valid desaturation data. The threshold value is 88% and the predetermined time is 300 seconds. Determining valid desaturation signals comprises eliminating artifacts. Determining valid desaturation signals comprises ignoring desaturation signals below the threshold value that do not reach a peak value.\nSome or all of the aspects of the invention described above may have some or all of the following advantages. The invention makes it impossible to automatically differentiate valid oxymetry data signals from invalid oxymetry data signals. The invention also automatically determines true desaturation signals from invalid desaturation signals.\nOther features, objects and advantages will become apparent from the following detailed description when read in connection with the accompanying drawing in which:"}
{"text": "Testing biological samples is typically a multi-step process performed in a laboratory by skilled technicians. The following workflow considers a DNA test for a human patient. The biological sample must be first taken from the donor using a swab or collection vial. The swab for wiping over buccal tissue is an open-celled foam which traps the donor cells within the foam cells. The foam is then immersed and agitated in a fluid to release the biological cells. A series of reagents* are added to (a) break the cells open and release DNA; (b) purify the released DNA and (c) mixed the DNA with reagents for amplification (pH, buffering, stabilizing agents, polymerase, primers, probes, beads, nucleotides, etc). These reagents are added either by hand pipetting or in an automated process for batch testing. At some point it is desirable to separate out unwanted particulate (food, cell debris, etc). This can be by a mechanical sieve or by centrifuge. The DNA may then be tested using known techniques such as Sanger sequencing, Sequencing by Synthesis, or real time PCR.\nThe entire process is quite complex and requires skilled persons using expensive equipment to provide a result, which typically takes hours to days to return to the person requesting the test.\nThe inventors appreciated that it is desirable to replace the above system of testing with a simple, inexpensive device not requiring specialist skills to operate or interpret. They have thus devised a device described below which can be hand operated by the lay user."}
{"text": "The present invention concerns generally a push-button control device of the type encountered particularly in electronic or electromechanical timepieces. The present invention also concerns a portable electronic instrument, such as a timepiece, including such a push-button control device.\nSuch push-button control devices are very widespread, particularly in the field of electronic or electromechanical timepieces. A typical example of such a control device is given in EP Patent No. 1 089 144. This document discloses in particular a multi-function push-button contact strip essentially formed of a base fixed to the bottom of a case in proximity to the push-button and an elastic blade forming a mobile contact element capable of being actuated by the push-button. This contact strip forms a first contact element of the control device, the elastic blade of the strip cooperating with a second contact element arranged on an electronic module for establishing an electric connection when the push-button is actuated.\nFrom the electrical and mechanical point of view, the contact is thus formed of two contact elements, one mobile (namely the aforementioned elastic blade capable of being actuated by the push-button) secured to the case by its base and the other fixed, which is arranged on the electronic module (typically a contact pad arranged on the printed circuit board of the electronic module). In the aforementioned EP Patent No. 1 089 144, the second contact element is for example formed of a bent conductive lug of the printed circuit. It will thus be understood that the two contact elements are arranged on different frames, namely the case (to which the push-button is typically secured) for the first contact element, and the printed circuit board, for the second contact element.\nOne drawback of the aforementioned typical construction lies in the fact that the printed circuit board constitutes an element whose dimensions and position in the case cannot be guaranteed with a high level of precision. This lack of precision leads to a loss of precision as to the arrangement and travel of the push-button, which consequently has repercussions on the quality of actuation and use of the push-button. In order to improve precision, it is thus necessary to find another solution or at least to improve the existing solutions.\nAnother drawback of the aforementioned typical construction also lies in the liability of the contact to wear. In fact, the fixed contact element with which the mobile contact element is brought into contact, is generally made in the form of a metallisation made on the printed circuit board or a part secured to the printed circuit. Following repeated mechanical contacts between the contact elements, this metallisation can undergo significant wear and consequently reduce the quality of the electric contact. It is thus also sought to reduce the liability of the push-button control device to wear."}
{"text": "People have long desired to be able to adjust the level of illumination of a light fixture so as to accommodate various lighting needs. For example, reading might require one level of illumination whereas background lighting might dictate a different level of illumination. To accommodate these differing lighting needs, various lighting arrangements have been developed which permit the level of illumination of the light fixture to be varied. It is believed that such existing lighting arrangements all rely solely on the use of dimmer switches to brighten or dim the intensity of the fixture's one or more light bulbs so as to adjust the level of illumination of the light fixture."}
{"text": "USB ports and cables allow interconnection of a variety of compatible electronic devices, such as desktop computers, automobile dashboard consoles and battery-powered portable devices such as laptop computers, tablets, mobile phones, e-readers and MP3 players. USB ports are accessed using standardized USB cable connections to provide serial communications between devices, as well as electrical power transfer for charging and operating battery-powered peripheral devices. USB compatible systems often include interface integrated circuits (ICs) mounted to an internal circuit board to interface USB data and power connections to host system circuitry such as power circuits and host processors. Dedicated USB charging devices are also available having multiple USB ports for charging various portable devices, which may include circuitry for fast charging certain peripheral devices. Many desktop and laptop computers include multiple USB ports for data transfer and/or peripheral device charging. USB power delivery (USB-PD) and Type-C (USB-C) specifications describe delivery of higher power over USB cables and connectors to provide a universal power plug for devices that may accommodate more than 5V charging power, for example, for fast or quick-charging capabilities. The USB-PD specification defines communications for negotiating voltage and current levels for power transfer from a source port to a sink port, where the power negotiation communications is independent from the normal USB communications.\nUSB-PD defines four kinds of USB compatible devices: Provider-Only, Provider/Consumer, Consumer/Provider, and Consumer-Only. Devices that are Provider-Only, Provider/Consumer, or Consumer/Providers may sometimes be in a Source role—meaning they are providing DC voltage on the Vbus wire for the far-end device to consume or sink. USB-C recognizes three kinds of USB devices: downward facing port (DFP), upward facing port (UFP), and dual-role port (DRP). In the absence of USB PD messaging, the DFP is the source of power and the UFP is sinking the power. USB-C cables and connectors include configuration channel (CC) lines for power configuration as well as for baseband communications. USB-PD specifications provide baseband communications using Biphase Mark Coding (BMC) for message exchange over a configuration channel (CC) wire or line of the USB cable. USB-C systems use a Type-C plug with two configuration channel lines CC1 and CC2. The USB-PD specification defines a half-duplex packet-based communication link between ports connected via a USB-PD cable and connectors to exchange information that enables the two ports to communicate and negotiate the voltage and current provided from a Source port to a Sink port. The ports can negotiate to switch roles (Source to Sink and vice versa). The BMC communications on the CC lines is independent from the normal USB communications that go through D+ and D− lines of the USB cable.\nThe CC line or lines may also be used for negotiating power transfer configurations of connected devices by way of analog signal levels. For example, up to 15 W of power can be delivered for USB Type-C cables without USB-PD messaging by controlling the DC voltage on the CC pin. The nominal voltage of the CC line is determined by pull up current from the DFP device (e.g., using a pull up resistor RP or a current source) and a pull down resistor RD (or pull down current source) from the UFP device. The CC line voltage value can thus vary from 0.3V to 2.4V in many instance due to combinations of the pull up and pull down levels. However, the baseband communication signals on the CC lines for typical BMC data packet exchange range from 0 to 1.1V. Consequently, the power supplies used for USB-C transmitters require pull down current to avoid having RP/RD combination charge up the baseband transmitter power supply node when a BMC output logic “1” is transmitted by the baseband transceiver. This is due to the CC line connection through the baseband transmitter to the supply that powers the transmitter. If inadequate pull down current is provided, and particularly if long duration transitions occur, the baseband transmit output can eventually be charged up and the transmit signals can go outside acceptable levels. For example, reverse current can be as high as 350 uA when a 47 k Ω pull up resistor and 3.3V supply are used. Accordingly, improved USB port controllers and techniques are desired to ensure robust baseband communications without excessive power consumption."}
{"text": "In the medical field, a medical instrument may be used as a blood circuit constituting an extracorporeal circulation circuit or the like in some cases. The medical instrument includes a tube body (hereinafter, referred to as a first tube body) which is flexible and a tube body (hereinafter, referred to as a second tube body) such as a connector, which has a rigidity higher than that of the first tube body (i.e., the second tube body is more rigid than the first tube body) and is fitted in the first tube body. When the medical instrument is used, blood circulates through a lumen of the first tube body and a lumen of the second tube body. Therefore, an inner peripheral surface of each tube body is provided with a coating layer having biocompatibility (i.e., a biocompatible coating is applied to the inner surface of the first and second tube bodies).\nThe coating layer is made of various types of materials which are available. For example, Japanese Patent Application Publication No. H04-152952 discloses a medical instrument that uses a synthetic polymer having biocompatibility and antithrombotic."}
{"text": "Over the years a variety of material dispensers have been developed including those directed at dispensing foamable material such as polyurethane foam which involves mixing certain chemicals together to form a polymeric product while at the same time generating gases such as carbon dioxide and water vapor. If those chemicals are selected so that they harden following the generation of the carbon dioxide and water vapor, they can be used to form “hardened” (e.g., a cushionable quality in a proper fully expanded state) polymer foams in which the mechanical foaming action is caused by the gaseous carbon dioxide and water vapor leaving the mixture.\nIn particular techniques, synthetic foams such as polyurethane foam are formed from liquid organic resins and polyisocyanates in a mixing chamber (e.g., a liquid form of isocyanate, which is often referenced in the industry as chemical “A”, and a multi-component liquid blend called polyurethane resin, which is often referenced in the industry as chemical “B”). The mixture can be dispensed into a receptacle, such as a package or a foam-in-place bag (see e.g., U.S. Pat. Nos. 4,674,268, 4,800,708 and 4,854,109), where it reacts to form a polyurethane foam.\nA particular problem associated with certain foams is that, once mixed, the organic resin and polyisocyanate generally react relatively rapidly so that their foam product tends to accumulate in all openings through which the material passes. Furthermore, some of the more useful polymers that form foamable compositions are adhesive. As a result, the foamable composition, which is often dispensed as a somewhat viscous liquid, tends to adhere to objects that it strikes and then harden in place. Many of these adhesive foamable compositions tenaciously stick to the contact surface making removal particularly difficult. Solvents are often utilized in an effort to remove the hardened foamable composition from surfaces not intended for contact, but even with solvents (particularly when considering the limitations on the type of solvents suited for worker contact or exposure) this can prove to be a difficult task. The undesirable adhesion can take place in the general region where chemicals A and B first come in contact (e.g., a dispenser mixing chamber) or an upstream location, as in individual injection ports, in light of the expansive quality of the mix, or downstream as in the outlet tip of the dispenser or, in actuality, anywhere in the vicinity of the dispensing device upon, for instance, a misaiming, misapplication or leak (e.g., a foam bag with leaking end or edge seals). For example, a “foam-up” in a foam-in-bag dispenser, where the mixed material is not properly confined within a receiving bag, can lead to foam hardening in every nook and cranny of the dispensing system making complete removal not reasonably attainable, particularly when considering the configuration of the prior art systems.\nBecause of this adhesion characteristic, steps have been taken in the prior art to attempt to preclude contact of chemicals A and B at non-desired locations as well as precluding the passage of mixed chemicals A/B from traveling to undesired areas or from dwelling in areas such as the discharge passageway for aiming the A/B chemical mixture. Examples of injection systems for such foamable compositions and their operation are described in U.S. Pat. Nos. 4,568,003 and 4,898,327, and incorporated herein by reference. As set forth in both of these patents, in a typical dispensing cartridge, the mixing chamber for the foam precursors is a cylindrical core having a bore that extends longitudinally there through. The core is typically formed from a fluorinated hydrocarbon polymer such as polytetrafluoroethylene (“PTFE” or “TFE”), fluorinated ethylene propylene (“FEP”) or perfluoroalkoxy (“PFA”). Polymers of this type are widely available from several companies, and one of the most familiar designations for such materials is “Teflon”, the trademark used by DuPont for such materials. For the sake of convenience and familiarity, such materials will be referred to herein as “Teflon”, although it will be understood that materials having the above and below described qualities are available from companies other than DuPont and can be used if otherwise appropriate.\nWhile features of the present invention are applicable to single component dispensing systems, the present invention is particularly suited for systems that have a plurality of openings (usually two) arranged in the core in communication with the bore for supplying mixing material such as organic resin and polyisocyanate to the bore, which acts as a mixing chamber. In a preferred embodiment of the invention, there is utilized a combination valving and purge rod positioned to slide in a close tolerance, “interference”, fit within the bore to control the flow of organic resin and polyisocyanate from the openings into the bore and the subsequent discharge of the foam from the cartridge.\nTeflon material and many of the related polymers have the ability to “cold flow” or “creep”. This cold flow distortion of the Teflon is both beneficial (e.g., allowing for the conformance of material about surfaces intended to be sealed off) and a cause of several problems, including the potential for the loss of the fit between the bore and the valving rod as well as the fit between the openings (e.g., ports) through which the separate precursors enter the bore for mixing and then dispensing. In many of the prior art systems utilizing Teflon, the Teflon core is fitted in the cartridge under a certain degree of compression in order to help prevent leaks in a manner in which a gasket is fitted under stress for the same purpose. This compression also encourages the Teflon to creep into any gaps or other openings that may be adjacent to it which can be either good or bad depending on the movement and what surface is being contacted or discontinued from contact in view of the cold flow.\nUnder these prior art systems, however, over time the sealing quality of the core is lost at least to some extent allowing for an initial build up of the hardenable material which can lead to a cycle of seal degradation and worsening build up of hardened material. This in turn can lead to a variety of problems including the partial blockage of chemical inlet ports so as to alter the desired flow mix and degrade the quality of foam produced. In other words, in typical injection cartridges the separate foam precursors enter the bore through separate entry ports. Polyurethane foam tends to build up at the area at which the precursor exits the port and enters the mixing chamber. Such buildups cause spraying in the output stream, and dispensing of the mixture in an improper ratio. The build up of hardened material can also lead to partial blockage of the dispenser's exit outlet causing a misaiming of the dispensed flow into contact with an undesirable surface (e.g., the operator or various nooks and crannies in the dispenser). Another source of improper foam output is found in a partially or completely blocked off dispenser outlet tip that, if occurs, can lead the foam spray in undesirable areas or system shutdown if the outlet becomes so blocked as to preclude output. A variety of prior art systems have been developed in an effort avoid tip blockage, particularly in automated systems, as in foam-in-bag systems, which impose additional requirements due to the typical high usage level and the less ready access to the tip as compared to a hand-held dispenser. The prior art systems include, for example, porous tips with solvent flush systems. However, over time these tips tend to load up with hardened foam and eventually become ineffective.\nThe build of hardened/adhesive material over time can lead to additional problems such as the valve rod and even a purge only rod, becoming so adhered within its region of reciprocal travel that either the driver mechanism is unable to move the rod (leading to an oft seen shut down signal generation in many common prior art systems) or a component along the drive train breaks off which is often the annular recessed valve rod engagement location relative to some prior art designs.\nThe above described dispensing device has utility in the packing industry such as hand held dispensers which can be used, for instance, to fill in cavities between an object being packed and a container (e.g., cardboard box) in which the object is positioned. Manufacturers who produce large quantities of a particular product also achieve efficiencies in utilizing automated dispensing devices which provide for automated packaging filling such as by controlled filling of a box conveyed past the dispenser (e.g., spraying into a box having a protective covering over the product), intermediate automated formation of molded foam bodies, or the automatic fabrication of foam filled bags, which can also either be preformed or placed in a desired location prior to full expansion of the foam whereupon the bag conforms in shape to the packed object as it expands out to its final shape.\nWith dispensing devices like the hand held and foam-in-bag dispensing apparatus described above, there is also a need to provide the chemical(s) (e.g., chemicals “A” and “B”) from their respective sources (typically a large container such as a 55 gallon container for each respective chemical) in the desired state (e.g., the desired flow rate, volume, pressure, and temperature). Thus, even with a brand new dispenser, there are additional requirements involved in attempting to achieve a desired foam product. Under the present state of the art a variety of pumping techniques have arisen which feature individual pumps designed for insertion into the chemical source containers coupled with a controller provided in an effort to maintain the desired flow rate characteristics through monitoring pump characteristics. The individual in “barrel” pumps typically feature a tachometer used in association with a controller attempting to maintain the desired flow rate of chemical to the dispenser by adjustment in pump output. The tachometers used in the prior art are relatively sensitive equipment and prone to breakdowns.\nIn an effort to address the injection of chemicals into the mixing chamber at the desired temperature(s) there has been developed heater systems positioned in the chemical conduits extending between the chemical supply and the dispenser, these heaters include temperature sensors (thermisters) and can be adjusted in an effort to achieve the desired temperature in the chemical leaving the feed line or conduit. Reference is made to, for example, U.S. Pat. Nos. 2,890,836 and 3,976,230, which references are incorporated by reference. These chemical conduit heater wires suffer from a variety of drawbacks such as (a) poor sensor (e.g., thermistors) responsiveness due to non head-on flow positioning of the sensor or difficulty in manipulating the sensor without breakage to be in the proper orientation, (b) difficulty in positioning the tip of the heater wire close enough to the dispenser to avoid cold shot formation and associated material stretch limitations in the heater wire conduit needed to avoid stretching and separation of the dispenser from the tip of the heater wire when the other “fixed” end originates from the pump control region, (c) increased pump weight and an increase in the length and cost associated with the leads extending from the heater wire tip to heater wire control and power source locations at the pump end, (d) an associated increase in electromagnetic interference (EMI) due to the longer “umbilical” cords and thermister leads, (e) poor thermister reliability in its heavy flex location within the interior of the heater wire, (f) difficulty in feeding heater elements within the outer protective chemical conduit, and (g) cost and production limitations in the overall heater wire and conduit length requiring relatively close positioning of the chemical driver source to the dispenser location.\nAs noted above, in the packaging industry, a variety of devices have been developed to automatically fabricate foam filled bags for use as protective inserts in packages. Some examples of these foam-in-bag fabrication devices can be seen in U.S. Pat. Nos. 5,376,219; 4,854,109; 4,983,007; 5,139,151; 5,575,435; 5,679,208; 5,727,370 and 6,311,740. In addition to the common occurrence of foam dispenser system lock up, cleaning downtime requirements, poor mix performance in prior art foam-in-bag systems, a dispenser system, featuring an apparatus for automatically fabricating foam filled bags, introduces some added complexity and operator problems. For example, an automated foam-in-bag system adds additional complexity relative to film supply, film tracking and tensioning, bag sealing/cutting, bag venting, film feed blockage. Thus, in addition to the variety of problems associated with the prior art attempts to provide chemicals to the dispenser in the proper rate, keeping the dispenser cartridge operational, and feeding film properly, the prior art foam-in-bag systems also represent a particular source of additional problems for the operators. These additional problems include, for example, attempting to understand and operate a highly complicated, multi-component assembly for feeding, sealing, tracking and/or supplying film to the bag formation area; high breakdown or misadjustment occurrence due to the number of components and complex arrangement of the components; high service requirements (also due in part to the number of components and high complexity of the arrangement in the components); poor quality bag formation, often associated with poor film tracking performance, difficulty in achieving proper bag seals and cuts, particularly when taking into consideration the degrading and contamination of heater wires due to, for example, foam build up and the inability to accurately monitor current heated wire temperature application, difficulty in formation and maintaining clear bag vent holes, as well as the inevitable foam contamination derivable from a number of sources such as the dispenser and/or bag leakage, and clean up requirements in general and when foam spillage occurs.\nAnother particularly problematic area associated with the prior art foam-in-bag system lies in the area of heated resistance wire replacement, both in regard to edge sealing and in regard to the cross-cutting sealing systems. In the prior art systems, there is often required delicate operator manipulation (see for example U.S. Pat. No. 5,376,219) with certain tools to achieve removal and reinsertion of broken, or worn, heated wires (which is a common occurrence in the thin heated resistance wires used in the industry to form the seals and cuts).\nIn addition, prior art systems suffer from other drawbacks, such as relatively slow bag formation and a slow throughput of completed bags which, in some systems, is partially due to a reverse feed requirement to break an upper, not-yet-completely formed bag from a completed bag adhered together by a bond formed by the earlier melted and presently cooled plastic material on the heated cross-cut wire.\nThe prior art mixing cartridge driver mechanisms for reciprocating valve rods has also shown in the field to be inadequate as they are subject to often breakdowns and often quickly become unable to achieve rod reciprocation after a minor build up of foam in the cartridge. An additional problem associated with the mixing chamber used on fixed dispenser embodiments such as a foam-in-bag dispenser is the difficulty in proper removal and mounting of a mixing module in the support housing. Prior art systems also suffer from hose and cable management (e.g., electronics, chemical supply and solvent supply) difficulties due to their becoming tangled and in a state of disarray so as to present obstacles to operators and potential equipment malfunctions due to cable or hose interference with moving components or the hoses/cables becoming disconnected and/or damaged.\nThe pump equipment of prior art systems are also prone to malfunction including the degrading of seals (e.g., isocyanate forms hardened crystals when exposed to air which can quickly degrade soft seals). The pumping systems currently used in the field are also subject to relatively rapid deterioration as they often operate at high rates during usage due to, for example, general inefficiency in driving the chemical from its source to the dispenser outlet. The common usage of in-barrel pump systems also introduces limitations in chemical source locations (e.g., typically a 20 foot range limitation for standard heater wire conduit and in barrel pump systems), which can make for difficulties in some operator facilities where it is required or preferred to have the chemical source located at a greater distance from the dispenser. The common usage of in-barrel pumps for prior-art dispenser systems also presents a requirement for multiple chemical sources to achieve the required one-to-one chemical source and pump combination, which in particularly problematic for operators running numerous dispenser systems.\nPrior art foam-in-bag systems, in presumably an effort to accurately dispense foam into the bag, locate the dispenser within the bag being formed (e.g., all dispenser components placed between the film left and right side edges and above the end seal of the bag). These prior art arrangements present problems from the stand point of the placement of the dispenser and its various components such as filters, chemical valving lines, and other components required for accessing a mixing module, all in the bag formation region. This positioning places those components in an area highly prone to chemical contact even with a properly functioning dispenser. Efforts have been made in the prior art to protect the dispenser through the use of covers, but these covers have shown to be highly ineffective in protecting the components. Once foam hardens on the components they are often made even more difficult to access when servicing is desired. Also, the non-smooth, multi-protrusion and edge presentment design of prior art foam dispensers, in addition to making cleaning impractical, have a tendency to create film tracking problems and/or require added guidance members to avoid film/dispenser contact.\nIn addition to the difficulty in achieving proper wire temperature levels in the chemical conduit heater wires, there has also been experienced difficulty in achieving proper end and edge sealing/cutting, and venting wire temperatures in prior art foam-in-bag systems. There is also associated with prior art systems problems in achieving proper positioning and in gaining access for servicing heater wires. The two most common prior art systems take different approaches with a first utilizing a rolling heater wire which presents added complexity in power supply as well as difficulty in removing and re-inserting heater wires. The second approach uses a non-rolling drag technique (e.g., U.S. Pat. No. 6,472,638) that, while being easy to remove and re-insert, has experienced difficulty in the field in maintaining a proper location of the exposed heater wire relative to the film being driven thereby, which is due in part to a tendency for the heated seal, wires becoming more and more embedded in the underlying support.\nFilm replenishment in the prior art systems has also proven to be difficult. Accessing prior art systems to remove the emptied roll and to replace it with a new role, which can be relatively heavy as in 25 lbs. or so, is only achieved with great difficulty due to the insertion location being in the rear, intermediate region of a typical foam-in-bag system design. This location is highly straining on the operator.\nMany prior art foam-in-bag systems and other automated dispending systems have shown in the field to have high service requirements due to, for example, breakdowns and rapid supply usage requirements (e.g., film, solvent, precursor chemicals, etc.). There is thus a great deal of servicing associated with prior art systems as in problem solving and in maintaining adequate supply levels. The prior art systems suffer from the problem of difficult and often non-adequate servicing which can be operator or service representative induced (e.g., failing to monitor own supply levels or anticipating level of usage or difficulty in responding timely to service requests which are often on an emergency or rush basis as any down time can be highly disruptive to an operator in timely meeting orders).\nAs can be seen there are numerous potential areas that can create problems in the field of dispensing."}
{"text": "1. Field of the Invention\nThe present invention relates generally to encoding and decoding for data transfer and, more particularly, to a method and apparatus for use in encoding and/or data by providing for a Hamming weight enhancement to the data.\n2. Related Art\nData is frequently stored and/or communicated in binary form, that is, as a sequence of binary digits or bits having zero and one values. When such data is communicated or simply transferred to a storage medium, such as occurs in the operation of a disk drive, for example, there is a substantial risk of transmission errors (e.g., a one-valued bit being received when a zero-valued bit was transmitted or vice-versa) caused by noise in the communication channel and other factors. These transmission errors are especially significant when data compression techniques are used to reduce the number of bits needed to communicate a particular message (e.g., in a disk drive or, more generally, in a communication channel), because a single erroneous bit in a compressed message can result in corruption of a larger amount of information in the message after the message is decompressed. Disk drives, of course, are well-known in the art, but additional information and a general circuit arrangement for one exemplary disk drive apparatus may be found, for example, in Wakamatsu U.S. Pat. No. 6,011,666, issued Jan. 4, 2000, the disclosure of which is hereby incorporated herein by reference as if set forth herein in its entirety.\nMany different approaches to overcoming this problem with transmission errors have been developed. FIG. 1 illustrates one example of a prior-art write-channel encoder apparatus for encoding information for communication via a channel (or for storage on any suitable media). As shown, an encoder 18 receives user data 20 and sequentially performs error correction coding or ECC (block 22), run length limit or RLL encoding (block 24), and precoding (block 26) on the user data 20 to thereby produce encoded data which is then written to any suitable media and/or any suitable communication channel (block 28).\nAn inverse process for decoding the information encoded by the encoder 18 of FIG. 1 is performed by a prior-art read-channel decoding apparatus such as the decoder 29 illustrated in FIG. 2. Initially, the decoder 29 reads or retrieves encoded information from a media source or a communication channel (block 30) and performs a maximum likelihood sequence detection (MLSD) on that encoded information (block 32). As will be appreciated by those of ordinary skill in the art, this MLSD detection compares the retrieved encoded information with a set of all possible information sequences and determines which of the possible sequences most likely represents the encoded information based on which of the possible information sequences is “closest” (e.g., by computing a predetermined error quantification metric for each possible sequence and selecting that sequence which has the smallest or minimum value for the computed error quantification metric).\nThe read-channel decoder 29 then performs reversed preceding on the detected data sequence (block 34) followed by RLL decoding (block 36) and ECC decoding (block 38) to thereby derive user data 40, which corresponds to the user data 20 encoded by the encoder 18 (FIG. 1).\nIn this approach, the RLL decoding step in the data recovery process undesirably causes small errors to be propagated into many symbols of error in the decoded message. As will be appreciated by those of ordinary skill in the art, the number of symbol errors that can be fixed by the error correction coding (ECC) is limited by design. Consequently, the RLL decoding used in this approach degrades the final error rate of the data recovery process.\nOne solution to this problem has been to reduce the size of the actual encoded word to minimize the maximum possible size for error propagation. In the most extreme case, only one symbol is encoded, and the rest of the symbols in the code word are left unencoded.\nBy way of example, where a communication scheme employs an eight-bit symbol size and an unencoded code word made up of 4 symbols {u1, u2, u3, u4}, the best error correction is obtained by simply encoding one of the symbols into a 9-bit code word. For example, if u1 is encoded to produce a 9-bit code word E1 a final code with low error-propagation can be constructed as follows:final code={u2, e1a, u3, e1b, u4}where ela is the first (most significant) 4 bits of the code word E1 and elb is the last (least significant) 5 bits of the code word E1. In such a scheme, following the decoding process, errors of small size (i.e., over 2 bits more than the smaller of ela or elb) can corrupt at most two symbols, which may occur when an error encompasses an unencoded symbol and one of the partially encoded symbols.\nThe advantage of this kind of simple encoding scheme is in the reduced error propagation after RLL decoding and thus the total number of symbols that can be corrupted by a single error event. As a consequence, this encoding scheme allows for an increase in the total number of error events the error correction coding (ECC) can correct. A significant disadvantage of this scheme is that the smaller encoded symbol makes the minimum Hamming weight of the total code word very small (i.e., equal to the Hamming weight of the single encoded symbol). As will be appreciated by those of ordinary skill in the art, the Hamming weight of a code word or other collection of symbols is defined as the number of non-zero symbols in the code word or collection of symbols in the Interleaved Non-Return to Zero Inverting (INRZI) domain. For binary signaling, the Hamming weight of a bit stream is the number of “1” bits in the bit stream. The read-channel device (i.e., the decoder) performs numerous functions, such as timing and gain recovery, before it passes sampled data to the MLSD detector. The purpose of RLL encoding is not only to remove quasi-catastrophic sequences that can degrade the performance of the MLSD, it is also used to introduce a guaranteed minimum number of signal transitions to allow the gain and timing loops to operate. Codes with high Hamming weight typically provide more signal transitions that can be used for the timing and gain loops. Codes with very low Hamming weight may not provide adequate timing and gain gradients for the channel to provide correct sampling phase and gain control.\nThis encoding scheme is least effective when every encoded symbol in the entire encoded data field has minimum Hamming weight and all the other unencoded symbols are zeros. The degradation in the timing and gain loops as a result of the channel operating on a sequence of low Hamming weight codes can be so large that it negates the benefit of the reduced error propagation made possible by this encoding scheme.\nTo mitigate this problem, the storage industry has used data-scrambling technique to minimize the probability of creating a large number of zeros in the user symbols. Unfortunately, this technique merely makes it more difficult for an end-user to locate a bad data sequence. There is still some finite probability that the bad data sequence can be created accidentally by the scrambling of some random data (e.g., random data created by running a string of zero-valued bits through the same randomizer used in the channel). This randomizer technique is therefore less than satisfactory.\nFIG. 3 illustrates an alternative embodiment of a prior-art write-channel encoder apparatus 41, which is similar to the encoder apparatus 18 of FIG. 1, except that the order in which the error correction coding (ECC) and the run length limit (RLL) encoding are performed is reversed in the encoding apparatus 41 of FIG. 3. In other words, as shown, run length limit encoding is performed on user data 42 (block 44), and error correction coding is subsequently performed (block 46). The result then undergoes precoding (block 48) and the encoded information is then written to any suitable media and/or a communication channel (block 50).\nSimilarly, FIG. 4 illustrates a read-channel decoding apparatus 51 which reads encoded information from a media or a communication channel (block 52) performs MLSD detection (block 54), reversed precoding (block 56), error correction code (ECC) decoding (block 58), and run length limit (RLL) decoding (block 60) to generate user data 62, which corresponds to the user data 42 processed by the encoding apparatus 41 (FIG. 3).\nPermuting or reversing the ECC and RLL encoders in the “write” process in this fashion serves to prevent or at least reduce error propagation. In such an encoding scheme, the ECC would correct on the data before RLL decoding, thereby completely avoiding the error propagation effect of RLL decoding. The output after ECC correction is assumed to be perfect user data, such that the RLL decoding following the ECC will not involve any error propagation. This method of recording and data recovery allows RLL code to be designed in such a way as to optimize for Hamming weights and thereby improve the performance of the timing and gain loops of the channel. Thus, the error propagation effect of the previously described encoding scheme are obviated.\nWhile this latter encoding scheme does provide at least a theoretical improvement in terms of reduced error propagation compared to the encoding scheme described above, it also has certain drawbacks due to compatibility issues that arise from the way in which storage devices that employ these encoding schemes are traditionally certified. Traditionally, technical customers of mass storage device often test the ECC capability of a given mass storage device before deciding to purchase the device. A common method of performing such a test involves so-called “read-long” and “write-long” processes.\nIn a read-long process, the drive or other device is instructed to read a certain data sector and return all the user symbols in that sector, plus the included ECC symbols, to the host. The host then runs a program that deliberately corrupts some number of bits in the received data (i.e., the user symbols and ECC symbols) and returns the corrupted data back to the drive using a write-long command. The write-long command instructs the drive to write all the long data back to the disk without additional ECC. In a subsequent normal read, the data is read back to the host and the host checks whether the ECC during a normal read has corrected the error that the host had deliberately introduced. This process is done repeatedly using long data corrupted at different location with different numbers of group of corrupted bits and different numbers of corrupted bits per group. These read-long and write-long instructions provide a means for the technical user to verify that the drive indeed has the specified error correction capability before deciding to purchase the drive.\nA problem with the permuted ECC/RLL scheme is that one corrupted symbol can generate a large number of encoded symbols that are corrupted, assuming that the RLL encoder used is the one that is optimized for the timing and gain loops of the channel. Thus, the conventional read-long and write-long test method will produce a result that appears to suggest that the ECC performance of the drive is vastly inferior to what is claimed by the manufacturer, and a technical customer may decline to purchase the drive because of this adverse test result. This, in turn, will cause the manufacturer of the drive to lose market share to other manufacturers that do not incorporate ECC/RLL permutation. For this reason, the permuted ECC/RLL scheme may not be readily accepted in the storage industry."}
{"text": "The invention relates generally to methods of non-destructive analysis. More particularly, the invention relates to an inspection method and device for quickly and non-destructively analyzing a container material for the presence of leaks.\nClosed packages are used to package a wide variety of products, including foods, condiments, and medicines, as well as other products such as batteries, and cosmetics. These packages function to contain and protect the products contained within them. They may also function as barriers to elements that might degrade, corrode, or otherwise compromise the quality of the product, including, for example, moisture, bacteria and other undesired microorganisms, and oxygen. The integrity of the packaging material is therefore of utmost importance, particularly when the contents are medicines. A liquid medicine contained within a package which has defects may slowly leak out undetected, resulting in a lower effective dosage. The medicine may become contaminated with a microbial agent, thereby compromising the health of the individual to whom the medicine is administered.\nVarious methods have been developed to detect leaks in packaging material. These methods include line scan camera surface inspection systems, LED systems, hydrocarbon leak detection systems, and SF6 detection systems. See, for example, U.S. Pat. Nos. 5,361,626; 5,499,529; 5,363,968; and 5,969,373. Other methods involve testing for flow of a gas such as CO2 or helium leaking from a package. Demorest et al. (1992) J. Plastic Film and Sheeting 8:109-123. These methods generally involve introduction of a gas into the package headspace by flushing the package and its contents prior to sealing; forcing the gas into the package after sealing with a pressurized bomb; or providing the gas as a normal component of the product or package. Typically, once the gas is introduced, a pressure differential is induced across the packaging material, and any gas leaks detected. Detection is generally accomplished by infrared detection device or by mass spectrometer. The detection limits of these methods are relatively high, however, and as a consequence, leaks may go undetected.\nThere is a need in the art for a rapid, sensitive, and effective inspection method which can analyze packaging material for defects before it is filled with a product.\nThe present invention provides methods and devices for detecting imperfections in a packaging material prior to filling with a product to be packaged. Packaging materials formed into a packaging unit are inspected or non-destructively analyzed by (1) applying a gas simultaneously with or immediately after forming the packaging unit; (2) detecting any gas which has passed through the packaging material of the packaging unit; and (3) analyzing the detected gas in a manner which makes it possible to quickly determine whether the packaging unit should xe2x80x9cpassxe2x80x9d inspection based on criteria such as amount of gas measured.\nIn some embodiments, the packaging unit is a blister, and is formed by mechanical deformation with a piston. The piston has a lumen extending through the piston, through which a gas is forced.\nDetection is carried out by a detection means. In some embodiments, the detection means is a thermal conductivity sensor. The detection means can detect leaks as small as about 1 xcexcm to about 0.1 xcexcm in the packaging material.\nThe detection means may be part of an inspection system, comprising a detection device, and a signal means for transmitting a signal to a filling means, which may be part of a manufacturing system. The signal means transmits either a positive or a negative signal to the filling means. A positive signal is generated in the event that no gas above detection limits is detected, and the filling means fills the packaging unit with product. A negative signal is generated in the event that gas is detected, and the filling means skips the defective packaging unit. Since containers with flaws are not filled with product, there is less product waste. Furthermore, since containers with flaws are not filled with product, the probability that the product will become contaminated with microorganisms, or will become chemically and/or physically altered before dispensing to a user is reduced.\nAn advantage provided by the inspection system of the present invention is that it can be integrated into a manufacturing system for the manufacture of containers, such as blister packs. Integration of manufacturing and inspection reduces inspection time.\nThese and other aspects, objects, advantages and features of the present invention will become apparent to those skilled in the art upon reading this disclosure and reviewing the drawings forming a part hereof."}
{"text": "This invention relates to methods for installing offshore tower structures.\nNumerous sea-related activities, such as oil exploration and recovery operations, for example, are conducted from offshore platform or tower structures. Towers have been employed which rest in an upright condition upon the water bed and are of extensive height, i.e., towers higher than 400 feet have been heretofore utilized.\nProblems of considerable magnitude have been experienced during the installation of mammouth offshore structures, giving rise to the proposal of various installation techniques, as demonstrated for example by the following U.S. Pat. Nos. 2,946,198, issued to Knapp on July 26, 1960; 3,633,369, issued to Lawrence on Jan. 11, 1972; 3,729,940, issued to Koehler on May 1, 1973; and French Pat. No. 1,444,839, issued May 31, 1966.\nIt has been proposed, for instance, to float an assembled tower in a horizontal position to an offshore worksite, upend the tower in the water and, thereafter, submerge the tower until its base rests upon the water bed.\nIt has also been proposed to float an assembled tower in upright fashion to the worksite and then gradually lower and immerse the tower onto the water bed.\nIt will be realized that massive pre-assembled towers, whether floated to a worksite in horizontal or upright positions, can be very difficult to support and maneuver, especially in rough seas. Moreover, the towers must be specially fabricated to withstand the high stress conditions occurring during transportation and immersion.\nIt has been suggested to install offshore structures by assembling component parts thereof at the worksite (see, for example, U.S. Pat. Nos. 2,534,480, issued to Shannon on Dec. 19, 1950 and 3,839,873, issued to Loire on Oct. 8, 1974.) Such techniques can be very time-consuming and may be hampered by unstable sea conditions, especially those involving structures which are highly susceptible to the effects of wave and wind action.\nIt is, therefore, an object of the present invention to eliminate or alleviate problems of the type previously discussed.\nIt is another object of the present invention to provide novel methods for installing offshore tower structures.\nIt is yet another object of the present invention to avoid subjecting offshore tower structures to high degrees of stress during installation.\nIt is a further object of the invention to support an immersed tower internally by means which can be subsequently removed.\nIt is still another object of the invention to provide novel methods and apparatus for the rapid on-site assembling of offshore tower components involving the use of removable buoyancy units which are floated within open-trussed legs of the tower and which are connectible to original leg segments and thereafter connectible to add-on leg segments of the tower to support the tower as it is gradually assembled and immersed."}
{"text": "The present invention relates generally to chopper stabilized amplifiers, and more particularly to circuitry and techniques for reducing the amount chip area occupied by capacitors used for Miller compensation and notch filtering.\nChopper stabilization is frequently used to improve offset and drift of amplifiers because it provides low in-band noise and avoids noise folding problems. A known technique for notch filtering can be used in chopper stabilized amplifiers to nearly eliminate the substantial output ripple voltage that otherwise occurs, and preserves the low offset, drift, and noise that are characteristic of the basic chopper stabilized amplifier topology. This known technique is described in the commonly assigned pending patent application Ser. No. 11/340,223 filed Jan. 26, 2006, entitled “Notch Filter for Ripple Reduction and Chopper Stabilized Amplifiers” by Rodney T. Burt and Joy Y. Zhang, incorporated herein by reference. However, increases bandwidth of a notch filtered chopper stabilized amplifier requires higher transconductance Gm in the chopping circuitry and accordingly requires larger compensation capacitors and larger capacitors in the notch filter. At higher bandwidths, the physical size of the compensation capacitors and notch filter capacitors causes them to occupy an unacceptably large amount of integrated circuit chip area, which unacceptably increases the cost of the chopper stabilized amplifier. This limits the usability of the otherwise highly desirable notch filtered chopper stabilized amplifier architecture in complex mixed signal systems in which it would be highly desirable to include a number of such amplifiers on a single integrated circuit chip.\n“Prior Art” FIG. 1 herein shows a somewhat simplified block diagram representation 1A of the same notch filtered chopper stabilized amplifier circuit shown in FIG. 3A of above mentioned pending application Ser. No. 11/340,223, including chopping circuitry, a notch filter and a feed forward stage. In Prior Art FIG. 1, feed-forward transconductance stage 5 is shown having a differential output between conductors 23A and 23B. Block 35 contains the input chopping switches 9, which are the same as switches 9-1,2,3,4 as shown in FIG. 3A of Ser. No. 11/340,223, coupled between the input voltage Vin and the inputs of a transconductance stage 2. Block 40 contains the output chopping switches 10, which are the same as switches 10-1,2,3,4 as shown in FIG. 3A of Ser. No. 11/340,223.\nIn Prior Art FIG. 1, chopping clock signal CHOPCLK on conductor 43 is applied to the control (CTL) inputs (not shown) of various input chopping switches 9 and various output chopping switches 10 and also is applied to the input of an inverter 41 which produces the logical complement of CHOPCLK and applies it to the CTL inputs (not shown) of various other input chopping switches 9 and the various other output chopping switches 10. CHOPCLK is the same as the “Phase1” clock signal shown in FIG. 3B of Ser. No. 11/340,223, and the logical complement of CHOPCLK is the same as the “Phase2” clock signal in FIG. 3B of co-pending Ser. No. 11/340,223. The input chopping switches 9 and output chopping switches 10 typically are implemented by means of individual MOS transistors or CMOS transmission gates. Similarly, filter clock signal FILTERCLK on conductor 51 is applied to the CTL inputs (not shown) of various switches 16 and 21 in notch filter 15 and also is applied to the input of an inverter 52 which produces the logical complement of FILTERCLK and applies it to the CTL inputs (not shown) of various switches 16 and 21 in notch filter 15, which can be the same as shown in FIG. 3A of Ser. No. 11/340,223.\nFILTERCLK is the same as “Phase3” in FIG. 3B of the above mentioned Ser. No. 11/340,223, and the logical complement of FILTERCLK is the same as “Phase4” in FIG. 3B of co-pending Ser. No. 11/340,223. Transconductance stage 3 (of transconductance gm2) and transconductance stage 4 (of transconductance gm3) in FIG. 3A of Ser. No. 11/340,223 are combined in block 3,4 in Prior Art FIG. 1 herein, wherein the two inputs IN1(+) and IN1(−) are the inputs of the “gm3” transconductance stage and the two inputs IN2(+) and IN2(−) are the inputs of the “gm2” transconductance stage. Output conductor 25 is connected to the output of the circuitry in block 3,4. Specifically, in block 3,4 the output of transconductance amplifier 4, referred to herein as “transconductance amplifier gm3”, in FIG. 3A of Ser. No. 11/340,223 is connected to Vout conductor 25. The (−) input of transconductance amplifier gm3 is connected by conductor 23 shown in FIG. 3A of Ser. No. 11/340,223 to the output of transconductance amplifier 5, to the output of transconductance amplifier 3, referred to herein as “transconductance amplifier gm2”, shown in FIG. 3A of Ser. No. 11/340,223 and also to one terminal of capacitor C1 shown in FIG. 3A of Ser. No. 11/340,223, the other terminal of which is connected to Vout conductor 25. The (−) and (+) inputs of transconductance amplifier gm2 are connected to conductors 22A and 22B, respectively.\nIn Prior Art FIG. 1, capacitors C0 and C1 are balancing capacitors, also referred to herein as “symmetrical counterpart capacitors”, to Miller compensation capacitors C2 and C3, respectively, and are provided to balance charge injection effects, to cancel a “zero” that would otherwise occur in the amplifier transfer characteristic, and also to provide good common mode rejection.\nFIG. 2 shows an implementation of the notch filter 15 shown in Prior Art FIG. 1, wherein the switches 16A-D and the switches 21A-D in FIG. 3A of Ser. No. 11/340,223 are implemented using N-channel MOS transistors.\nThus, there is an unmet need for circuitry in a chopper stabilized amplifier which substantially reduces the amount of integrated circuit chip area required for the various capacitors associated with notch filtering and Miller compensation within the chopper stabilized amplifier."}
{"text": "Many states have adopted Renewable Portfolio Standards (RPS) that require electricity supply companies to increase energy production that is attributed to renewable energy sources. The federal government may soon implement a renewable electricity standard (RES) that would be similar to the “renewables obligation” imposed in the United Kingdom. These standards place an obligation on electricity supply companies to produce a specified fraction of their electricity from renewable energy sources, such as wind, solar, hydroelectric, geothermal, biofuels, and biomass.\n“Biomass” refers to renewable organic materials such as wood, forestry waste, energy crops, municipal waste, plant materials, or agricultural waste. Biomass often contains about 10 to about 50 weight percent moisture. The trapped moisture cannot be used as fuel and increases costs associated with transportation of the biomass. Thus biomass is a low grade, high cost fuel that cannot compete economically with the fuel most commonly used to generate electricity—coal. Further, biomass has a low bulk density, is very hydrophilic, is seasonal, is variable, and has a limited shelf life.\n“Torrefaction” refers to the processing of biomass at temperatures between about 200° C. to about 350° C. (400°-660° F.) at atmospheric pressure wherein water and light volatile organic chemicals associated with the raw biomass material (i.e., “feed stock”) are vaporized. In addition, during the torrefaction process, molecules of biopolymers (cellulose, hemicelluloses and lignin) contained in the biomass decompose. After torrefaction, the biomass is a solid, dry, blackened material that is often referred to as “torrefied biomass” or “biocoal” that is easier to grind, which allows it to be used in coal burning power plants. Further, the torrefied biomass possesses a lower oxygen content, has a significantly reduced moisture content (less than about 3%), and has higher fixed carbon levels, which is directly proportional to heating value.\nFluid bed reactors are commonly used to carry out multiphase reactions. In this type of reactor, gas or liquid is passed through a granular solid material at high enough velocity to suspend the solid and cause it to behave as though it were a fluid. This process, known as “fluidization” imparts many important advantages to the reactor. As a result, the fluidized bed reactor is now used in many industrial applications, such as coal drying. Commonly coal drying is performed in an inert gas, i.e., oxygen-free environment. Drying coal in a non-oxidizing environment requires external heat sources to maintain the temperature of the reactor. However, coal has been dried in an oxidizing environment where the heat used to support the process is at least partially drawn from the burning coal. The temperature of the fluid bed reactor used to dry and otherwise process the coal is controlled by balancing the rate at which the coal is fed into the reactor against the amount of heat generated by the combustion process. Drying of coal increases the heating value of low rank coals, reduces the particle size of the feed stock, and partially decarboxylizes and desulfurizes the coal. After the coal is dried, it must be rehydrated to raise the moisture content up to about 5-9% to reduce its spontaneous combustion characteristics so that it is similar to native coal.\nThe table provided below illustrates the differences between raw coal and processed coal. One of skill in the art will appreciate that processed coal possesses a higher fixed carbon and heating values correspond to raw coal and the moisture content is drastically reduced.\nRaw CoalProduct 1Product 2Product 2Proximate Analysis:Moisture20.16%8.00%8.00%8.00%Ash8.16%7.93%8.69%8.67%Volatile Matter31.70%35.33%34.90%35.05%Fixed Carbon39.98%48.74%48.42%42.48%Ultimate Analysis:Moisture20.16%8.00%8.00%8.00%Hydrogen2.87%3.32%3.19%3.14%Carbon55.50%63.15%62.65%62.74%Nitrogen0.75%0.99%1.12%0.81%Sulfur0.77%0.52%0.54%0.48%Oxygen11.79%16.09%15.82%16.16%Ash8.16%7.93%8.69%8.67%Heating Value,9,44410,46010,31510,165Btu/lb"}
{"text": "1. Field of the Invention\nThis invention relates to a compact, lightweight device to produce high-power acoustic signals of relatively low frequencies in water.\n2. Description of the Related Art\nAcoustic devices used to produce high-power, low-frequency sound in the ocean are characterized by their large size and weight. It is desired to be able to generate hundreds of watts of omnidirectional acoustic power at frequencies below 1000 Hz with a device that can be deployed from an aircraft, surface vessel, or submarine. Although conventional low-frequency sources can usually operate over a broad band of frequencies, there exists a definite desire for a compact device with more limited bandwidth which could be assembled in modular configurations of units with different frequencies to meet frequency spectrum requirements.\nIn order to create high-power acoustic tones in water at low frequencies (below 1000 Hz), a device must produce large volume displacements. The volume displacement is the integral of the normal displacement of a radiating area, taken over that area. Therefore, an acoustic source must either have a large radiating area, a large displacement, or a combination of both.\nIn order to be effective for such applications as anti-submarine warfare and undersea geological surveying, a transducer must be operable at a wide variety of ocean depths, including great depths of several thousand feet. It is also desired to have such a transducer with an output that varies linearly with the transducer input."}
{"text": "The present invention relates to systems for correlating patterns with images.\nIn image correlation, a desired pattern (e.g., a desired object to be picked up by a robot arm) is compared with an image from some source (e.g., a camera in an automated assembly line). Because of the large amount of data in the form of pixels in an image to be processed for image correlation, the employment of parallel processing for this comparison is desirable.\nThere exists a great commercial need for faster and more flexible computer vision systems. For many applications within automated assembly lines, existing computer vision systems are too slow, too limited, or both. Highly parallel VLSI (very large scale integration) based vision systems have the potential for achieving the speed and complexity advances required. Unfortunately, most computational systems do not do well in terms of system performance when applied in parallel. In the general case this is usually due to the lack of sufficient realizable parallelism in the problem to be solved, but vision is a high parallel activity. Existing parallel vision systems are able to do the initial massive parallel computation of the raw image data, but are ineffective in combining the results of the low level computations, analyzing the results, and passing the results of the analysis on to higher processing stages. Indeed complex vision processing can produce much more intermediate data than the amount of data in the initial image. If the architecture of the VLSI vision system is not designed correctly, the fundamental processing bottleneck in the system can be the passing of the intermediate data to the chips where it is needed.\nThe Fairchild Bit-plane Area Correlator (BAC) is an example of a circuit which does the actual comparison of image pixels to pattern pixels. The BAC was designed with the above problems in mind, and processes most of the image data directly on one chip, passing only limited higher level results back outside the chip. The BAC is a high speed engine for computing correlation style match functions between two images passed to it and was designed with parallel deployment in mind. To utilize the BAC in a pattern recognition system, a method of rapidly supplying a number of BACs in parallel with input images and stored patterns is needed.\nMany image processing applications involve simply looking for the presence (or absence) of a known pattern and determining its location. Examples include locating parts on a conveyor belt, verifying the presence of parts during assembly operations, locating bonding pads and alignment marks on semiconductor chips, locating landmarks or targets in aerial imagery, and optical character recognition. Many video signal processing applications involve the computationally related process of convolving a two-dimensional signal array with a mask (e.g., to remove noise or enhance edges).\nBoth above classes of application require high-speed correlation (or matching) of a mask with an image. In addition, for the image processing examples, it is also necessary to determine the location with the best match. Traditionally, such functions have been implemented in special purpose boards using off-the-shelf components, resulting in large, expensive systems.\nAn example of how a correlation might be done is shown in FIGS. 7 and 7A. A rectangular image 210 in an image field 212 is captured by a video camera or other means, preprocessed to remove noise or other distortion or enhance edges, etc., and is then presented to a correlator. The image field 212 is represented as a series of binary digits as shown in field 214. Here the zeroes would represent a white background while the ones would represent the dark object.\nIn the correlation process, image 210 is compared with a series of masks 216, 218, 220, 222, and 224, which are stored in memory as an array of zeroes and ones. The masks have stored different orientations of an ideal image about one axis. The remaining two axes can be matched by moving the image horizontally and vertically, as shown in FIG. 7A. In FIG. 7A, image field 212 with image 210 is being matched against mask 220. As can be seen, the image is first shifted to the right as shown by fields 226, 228, and 230, and then is shifted upward until a best match is obtained. A simple way of determining a match is to compare each digital bit in image field 212 to each digital bit in mask 210 and produce a count of the number of bits which match. This count can then be compared with a threshold value which is chosen for a count which is close enough to indicate a match of the image with the mask. In one alternate method, a small portion of the ideal image could be stored as a mask (such as a 64.times.64 mask for a 128.times.128 image). Portions of the image can then be compared with the mask to determine where a match occurs.\nEach line of the digital representation shown in field 214 is referred to as a scan line and the entire field 214 is referred to as a frame. The traditional method of shifting the image as shown in FIG. 7A is to input the digits of a scan line into a series of shift registers and shift the digits to the right sequentially and compare the digits to the mask at each shift position. Such a correlation method is shown in U.S. Pat. No. 4,200,861 to Hubach, et al. As can be seen, this requirement of shifting the bits of an image slows down the processing time due to the requirement for doing a comparison at each shift position.\nThe processing time required is further increased if one attempts to do gray scale matching. For gray scale matching, rather than each pixel of the image field 212 being either a digital zero or one, each pixel will be a binary number ranging, for instance, from zero to sixteen. Zero could represent white, with sixteen representing black, and the numbers in between representing different shades of grey. Each pixel is then represented by a four bit binary code as illustrated in FIG. 8. The correlation can then be done by comparing bits in a bit plane which consists of a corresponding bit in each of the pixels. The bit planes compared would be first the most significant bit in every pixel, with subsequent comparisons being made of the lesser significant bits. As can be seen, the number of shifts through a shift register required for a correlation is increased by severalfold, with a corresponding increase in processing time."}
{"text": "Reversible and returnable envelope blanks are known wherein tabs or flaps are folded in one way to form an envelope sent by the sender to the receiver and then re-folded in another way to form a return envelope for re-mailing back to the sender. For example, U.S. Pat. No. 4,210,250 describes a generally cross-shaped envelope blank having a sending side and a returning side and two pairs of opposed flaps extending from a central panel. There are a total of five deposits of adhesive on various of the flaps, two of which are of an always-sticky gum-type releasable adhesive which once used and separated losed its adhesive properties, and three of which are of a moisture-activatable adhesive which become sticky only when moistened and which after drying forms a permanent bond. Such moisture-activatable adhesives have the disadvantage of being susceptible to unintentional activation from exposure to water or humidity. Also, release paper is typically used on the gum-type releasable pressure sensitive adhesives. Moreover, following the correct procedure in folding the flaps of this prior art blank together at the first mailing and then again upon re-mailing is by no means uncomplicated and if it is not done properly the first time the envelope must be discarded because none of the adhesive deposits can be used twice.\nOther returnable envelopes disclosed in the prior art include that of U.S. Pat. No. 2,829,818 which involves stapling parts of the blank together when it is readied for re-mailing. Some prior art reversible envelope designs include tear strips such as those disclosed in U.S. Pat. Nos. 4,669,652 and 4,917,287.\nIt is a principal object of the present invention to provide a blank for forming a reversible and returnable envelope which avoids the use of staples and tear strips, avoids moisture-activated adhesive susceptible to unintended activation and avoids the use of release paper usually required for pressure-sensitive adhesives. A further object is to permit the envelope to be folded and re-folded such that if the folding is done mistakenly the flaps can be peeled apart and done again correctly, which is not possible with adhesives which can only be used once. The invention further includes among its purposes an economical use of both paper and adhesive, avoidance of the danger of tearing of the sheet material of the blank and avoidance of exposure of any adhesive not in use during sending or returning."}
{"text": "The field of the present invention is building construction hardware for mounting components such as utilities to steel web joists.\nBuilding construction frequently uses steel web joists. Such devices employ two beams. The beams are mutually parallel and spaced apart to create a large moment of inertia for the joist. Cord elements extend at angles between the two so the beams will act as one in bending. Each beam includes a set of two elongate angle elements which are also mutually parallel. The cord elements extend for anchoring to between the elongate angle elements which form a cord space to receive the elements. Each angle element includes two legs extending at a right angle to one another, a first leg which is parallel with the first leg of the other angle element to define the cord space and a second leg extending in opposite directions from the second leg of the other angle element to generally define laterally extending flanges lying in a plane. Above the angle element of the upper beam, a wood nailer may be employed. The wood nailer sits on the plane defined by the laterally extending flanges and is typically approximately the same width as the beam defined by the two flanges.\nRight angle brackets having holes centrally extending through each of the two legs of the bracket are well known. Such brackets are used in both wood and steel construction. However, steel construction often requires the hanging of various construction components such as pipes from the basic steel structure. The steel structure in such construction typically is not conveniently or economically modified through drilling, welding or the like to provide attachment for such components. Consequently, it is advantageous to provide clamping mechanisms, which can, without modifying the beam, rigidly support a mounting for the beam. Such clamping mechanisms frequently are adjustable to accommodate the variations in such beams and avoid the need for stocking excessive parts. An adjustable clamping mechanism providing a mounting is illustrated in U.S. patent application Ser. No. 09/389,771, filed Sep. 2, 1999, for Hanger For Mounting Two I-Beams, the disclosure of which is incorporated herein by reference.\nThe present invention is directed to a seismic adapter for attachment to the beam of a steel web joist including a wood nailer and a method of mounting of such an adapter.\nIn a first separate aspect of the present invention, a seismic adapter includes two brackets positionable to either side of the wood nailer and retained by a shaft extending through the wood nailer. The brackets each include a mounting section to receive the shaft, a base section to extend beneath the beam of the steel web joist and an engagement section to engage the underside of the laterally extending flange of the steel web joist.\nIn a second separate aspect of the present invention, a seismic mounting includes two brackets positionable to either side of the wood nailer and retained by a shaft extending through the wood nailer. The brackets each include a mounting section with a hole to receive the shaft, a base section to extend beneath the beam of the steel web joist and an engagement section to engage the underside of the laterally extending flange of the steel web joist. The shaft extends through a passage in the wood nailer at a center distance above the undersurface of the laterally extending flanges that is greater than the distance from the axis of the holes located on the mounting section of each bracket to the edges of the engagement sections of each bracket.\nIn a third separate aspect of the present invention, a method of seismic attachment employs the placement of brackets to either side of the wood nailer such that base sections extend inwardly beneath the laterally extending flanges and engagement sections rise to meet the underside of the flanges. A passage extends transversely through the wood nailer and a shaft through the passageway provides for drawing the mounting sections against the wood nailer. The passage is drilled such that the center distance above the undersurface of the laterally extending flanges of the steel web joist is greater than the distance from the axis of the holes through the mounting section to the edges of the engagement sections. In this way, the edges of the engagement sections are drawn hard against the underside of the laterally extending flanges of the steel web joist beam.\nIn a fourth separate aspect of the present invention, the foregoing separate aspects are contemplated to be employed in combination.\nAccordingly, it is an object of the present invention to provide an improved seismic adapter and a method of seismic attachment to a steel web joist. Other and further objects and advantages will appear hereinafter."}
{"text": "MEMS devices comprise a relatively new technology that combines semiconductors with very small mechanical devices. MEMS devices are micro-machined sensors, actuators, and other structures that are formed by the addition, subtraction, modification, and patterning of materials using techniques originally developed for the integrated circuit industry. MEMS devices are used in a variety of applications, such as in sensors for motion controllers, inkjet printers, airbags, microphones, and gyroscopes. MEMS devices are increasingly used in a variety of applications such as mobile phones, automobiles, global positioning systems (GPS), video games, consumer electronics, automotive safety, and medical technology. Many potential and current applications require integration of MEMS devices with other types of chips or functionality. For example, MEMS devices may be integrated with bipolar, CMOS logic, or other peripheral devices such as trench or MIM capacitors.\nManufacturing MEMS devices is challenging in many aspects. Fabricating small moving parts of MEMS devices with lithography processes used in semiconductor technology has limitations. For example, lithography systems and processes restrict the minimum gap between moving and stationary parts of MEMS devices. Further, for reducing cost of these devices, it is imperative to minimize manufacturing costs.\nThus, what is needed in the art are cost effective structures for MEMS devices and methods of manufacture thereof."}
{"text": "Balloon catheters are increasingly being employed to conduct neurological procedures in patients. However, the design parameters for balloon catheters intended for use in neurological procedures are significantly different than the design parameters for balloon catheters used in non-neurological procedures such as cardiological procedures. For example, the width of the circulatory system within the neuroanatomy is significantly smaller and more tortuous than the circulatory system in other parts of the body. In order to access the smaller and more tortuous regions of the neuroanatomy, it is necessary to minimize the outer diameter of the balloon catheter while simultaneously maintaining the pushability and trackability of the catheter.\nIn order to minimize the outer diameter, current balloon catheters intended for neurological procedures employ a non-reinforced, single lumen, over-the-wire design. Accordingly, these balloon catheters are prone to several problems. First, the non-reinforced lumen is susceptible to ovalizing and/or kinking which, in turn, hinders advancement of the catheter over the guidewire, as well as deflation of the balloon. Second, the single lumen is in communication with the arterial blood flow. As the guidewire and balloon catheter are manipulated through the circulatory system, blood is withdrawn into the single lumen of the balloon catheter. Blood may thereby enter the balloon during inflation and cause (1) poor imaging of the balloon, for example, poor fluoroscopic imaging; (2) poor passage of the balloon through the circulatory system due to the premature inflation of the balloon; and (3) poor deflation of the balloon due to blood coagulation in the balloon inflation/deflation port. An additional disadvantage of single lumen balloon catheters is that the interference fit of the guidewire and inflation seal of the balloon may result in removal or peeling of the hydrophilic coating of the guidewire.\nIn order to minimize the outer diameter, current balloon catheters intended for neurological procedures are also typically designed to work with only a narrow gauge guidewire that is supplied by a manufacturer along with the balloon catheter. The current balloon catheters employ guidewires having diameters in the range of 0.010 to 0.012 inches. These relatively narrow guidewires are soft and, therefore, are very difficult to maneuver through the small, tortuous neuroanatomy.\nWhat is needed in the field is a balloon catheter that is operable to use with larger gauge guidewires; resists ovalizing and kinking of the inflation and guidewire lumen(s); and deploys with improved pushability and trackability."}
{"text": "This invention relates to fixation of objects to bone.\nA variety of techniques are available for affixing objects, such as soft tissue, to bone. The oldest technique for affixing soft tissue to bone is to pass a thread through the bone and sew the tissue down to the bone. Many sizes, shapes and types of suture and needles are available to accomplish this task. Today, this method is still used for repair of tendons and ligaments in the hands of older osteoarthritic patients, although passing sutures through bone is generally difficult and tedious.\nVenable et al. (1937), Annals Surg., Vol. 105, pp. 917-38, describes reactions of bone to metals, and recommends using Vitalium, an iron free and body fluid resistant alloy, for internal fixation devices, and particularly for screws. Prior to this publication, soft tissue repairs generally had been accomplished with string (suture), silver wires, or ivory pegs. Following publication of Venable et al., screws came into common surgical use, at first for repairing bone fractures, and later for attaching soft tissue to bone.\nVitalium staples were described for use in soft tissue repairs about 1940. Later many different types and styles of staple were suggested, including, for example, staples having four legs, or barbs in the legs, or various configurations of crossmember, including dentitions.\nMetal implants, including implants made of stainless steel, are subject in time to corrosion and consequent loss of structure. Moreover, the presence of metal in an anatomical site can interfere with certain imaging diagnostic or therapeutic treatments near the site, such as magnetic resonance imaging; and where the use of such imaging is indicated, any metal implants may first have to be surgically removed. Patient sensitivity to free nickel ions in stainless steel implants has fueled a growing controversy regarding the use of materials containing high quantities of nickel, including nickel-titanium alloys such as Nitinol.\nGenerally, the tissues react to metal screws and staples as foreign bodies, and these objects can be expected to loosen in time as a result of rejection of them by the tissues. A loosened screw or staple can be expelled from the bone, and can lodge in a joint, where it can cause damage to articulating surfaces. In a significant proportion of cases, where the screw or staple has been expelled or has loosened, it must be removed in a subsequent surgical procedure.\nSince the development of the stainless steel screw and staple, many small improvements have been made for fixing soft tissue to bone. In one improvement, described for example in Daniel et al., Chapter 8, In: Jackson et al. (1987), The Anterior Cruciate Deficient Knee, C. V. Mosby, pp. 114-126, a circular plastic washer with spikes on its undersurface is installed beneath the screw head to provide fixation without crushing the tissue, and toothed washer devices have been commercialized by, for example, A. O. Synthes, Switzerland. This method was until recently widely accepted as the best method of fixating soft tissue to bone, except in shoulder repair where sutures passed through holes drilled through the glenoid margin and through the edge of the glenoid capsule continue to be used for approximating the capsule to the glenohumeral rim, generally as described in Bankart (1983), British Jour. Surg., Vol. 26, pp. 23-39. Rowe et al. (1984), Jour. Bone Joint Surg., Vol. 66A, pp. 159-68, for example, describes using a Bankart procedure and in all instances avoiding the use of metal implants such as staples and screws in the vicinity of the shoulder joint.\nNecrosis of the soft tissue can result if the tissue is too tightly clamped by screws or staples, and several attempts have been made recently to improve the soft tissue fixation by screws and staples to overcome this problem, such as by using toothed washers, as described above, or by using stand-off devices to prevent crushing the soft tissue. On the other hand, if the soft tissue is too loosely fixed to the bone, the holding power of the tissue attachment is inadequate to facilitate effective soft tissue reattachment to bone. Because it is almost impossible to adjust the compression exerted by screws and staples on soft tissue, these devices are not fully satisfactory for soft tissue repair. The surgeon's fondness for suture in soft tissue repairs has never diminished, owing primarily to the fact that, by setting the tension of the suture, the surgeon can fix the soft tissue to the bone as tightly as is appropriate in the particular case and according to the surgeon's practice.\nSomers et al. (1985), U.S. Pat. No. 4,632,100, describes a cylindrical suture anchor having a drill formed at one end and flights of screw threads at the other end. The device combines a drill point for penetrating the hard outer cortical bone with a screw for fixing the device into the hard bone, providing for drilling the bone and installing the anchor in one operation. A Somers et al. device, marketed as STATAK.TM. by Zimmer, Inc. (Warsaw, Ind.), is screwed into a 4.5 mm diameter bone hole, and is intended to be countersunk into the cancellous bone to a depth up to 18 mm, and so it is not ideal for use in smaller joints. The relatively large size and the comparatively high cost of the device, and its requirement for a large hole in the bone, may be deterrent to its use, and it has not been well received.\nGoble et al. (1986), U.S. Pat. No. 4,738,255 describes a suture rivet that can be inserted into and locked in place in a performed hole in bone. The initial pilot hole is cut using a drill having scissoring blades, so that the resulting hole is flared or skirted, having a greater diameter deeper within the bone than at the entry to the bone surface. The anchor itself includes a rivet and a slotted ring. The slotted ring, which passes the smaller entry hole, is inserted within the bone, and then the rivet is passed into the hole and through the ring, fracturing the ring at the point where it is slotted and flaring the ring so that it cannot pass out through the entry hole, locking the assembly into the bone hole. The Goble et al. apparatus is relatively complicated, and the scissoring device can be unreliable in use.\nGatturna et al. (1987), U.S. Pat. No. 4,898,156, describes a suture anchor assembly consisting of a titanium body affixed to a nickel titanium arc of wire. The nickel titanium arc is of pseudoelastic nitinol, which can be strained such that the arc will straighten completely, allowing the assembly to be placed into a predrilled bone hole. The titanium body has a drilled cross hole which allows a suture to be attached to the body prior to insertion in the bone hole. The installation tool consists of coaxially disposed cannulae, which constrain the nitinol arc prior to insertion.\nNicholson et al. (1987), U.S. Pat. No. 4,899,743, describes a suture anchor installation tool that holds the Gatturna et al. suture anchor so that the nitinol arc is unconstrained, essentially hanging outside the installation tool. Prior to insertion the assembly resembles a one barbed fish hook in the end of a small spear. A hole is predrilled in the bone, and the anchor assembly is gently pushed into the predrilled hole, allowing the nickel titanium arc to slide down the side of the hole. A Gatturna et al. anchor and Nicholson et al. installation tool, commercialized by Mitek Surgical Products, Inc., beginning in late 1989, have become a preferred method of anchoring soft tissue to bone; many were sold in the first year following FDA approval. The device can be effective in aiding in soft tissue reattachment, and has an advantage in requiring a relatively narrow hole in the bone.\nHowever, the Gatturna et al. anchor is set by applying traction to the suture, and so the device does not lend itself to use as a rivet, in which no suture would be required to anchor the tissue. Moreover, use of the device is contraindicated in very soft, cancellous bone, as is typical of bone in many joints, so the device is not optimal for such procedures as rotator cuff repair, osteoarthritis joint reconstruction, and the like. When the Gatturna et al. anchor is deployed in an arthroscopic application using the Nicholson et al. installation tool, the protruding barb can engage sensitive tissues during placement and cause damage to the tissues.\nHayhurst (1988), U.S. Pat. No. 4,741,330, and international patent Publication No. WO 89/10096 describes a suture anchor including a generally bullet-shaped resilient plastic member having a rounded convex base from which legs extend. The legs, which are provided with outward-pointing barbs on their outer surface, diverge outwardly when the member is in a relaxed state. The member is compressed and inserted into a predrilled bone hole, and then allowed to relax, so that its resilience is said to urge the legs outward against the bone hole wall. The anchor supposed to be set by applying tension to the suture, causing the edges of the legs and the surface barbs to dig into the bone. A similar device, marketed by Acufex Microsurgical, is driven into the bone hole apparently into the cancellous bone, and is then set by pressing a spreader downward between the legs while pulling upward on the suture to force open the legs.\nBays et al. (1990), U.S. Pat. No. 4,924,865, describes a bioabsorbable tack for joining severed or torn soft tissues, such as cartilage. The stem of the generally T-shaped tack is hollow and covered with barbs on its outer surface. The tack is implanted by passing the stem through a hollow cylindrical applicator and passing a needle through the axial bore in the stem, then driving the needle, tack and applicator through the tissues to the desired depth, and then withdrawing the needle and applicator, leaving the tack in place."}
{"text": "The present invention relates generally to electric heaters and more specifically to positive temperature coefficient (PTC) electric heaters.\nTwo well-known types of electric heaters are the resistance-wire heater and the positive temperature coefficient (PTC) heater. In a resistance-wire heater, heat is produced by passing current through an electric wire. In a PTC heater, heat is produced by passing current through a plurality of electrically interconnected PTC (ceramic)pellets. The primary drawback of the resistance-wire heater is its propensity for overheating, which can often result in a fire if there are flammable items within the heater\"\"s immediate vicinity. More specifically, the heat source in a resistance-wire heater can often reach temperatures as high as 1600 degrees fahrenheit, well above the limit of combustible materials. Without installing a an effective safety mechanism in the resistance-wire heaters, the possibility of overheating and fire is present.\nIn contrast to the resistance-wire type heaters, the possibility of overheating and fire in a PTC heater is not present. This is a result of the self-limiting characteristic of the PTC element\"\"s temperature. Due to the physical chemistry of the doped ceramic material of which the PTC elements are made in a PTC heater, the PTC elements effectively shut off at a relatively low temperature. Typically, once a PTC element\"\"s temperature rises to approximately 340 degrees fahrenheit, the resistance of the element rises to a point to completely eliminate the flow of current through it, thereby preventing any further rising of the FTC element\"\"s temperature. This physical characteristic of a ceramic FTC heater element prevents overheating and eliminates the possibility of fire.\nIn addition, FTC heaters are more durable than resistance-wire heaters. A limitation of PTC heaters, however, is that the heat produced is directed over a relatively small area. This is due primarily to the size and shape of the heater assembly containing the ceramic pellets.\nSome pertinent patents include U.S. Pat. Nos. 5,245,692, 5,197,112, 4,965,432, 4,737,616, 4,703,153, 4,518,847, and 4,459,466.\nIt is an object of the present invention to provide a new and improved PTC heater.\nIt is another object of the present invention to provide a PTC heater which is small in volume and which can generate a large amount of heat.\nIt is still another object of the present invention to provide a PTC heater which disperses the heat emitted over a large area.\nIt is yet another object of the present invention to provide a PTC heater which can be mass produced, has a minimal number of parts, and can be easily assembled.\nIn furtherance of the objects broadly set forth above, a positive temperature coefficient (PTC) heater is provided which comprises a housing having a front and a rear, the front and rear each having at least one opening therein, a fan mounted within the housing for producing a flow of air in from the opening in the rear and out through the opening in the front, a pair of positive temperature coefficient (PTC) heater assemblies disposed in the housing in front of the fan in side by side relationship and angled relative to each other, and a plenum mounted in the housing for directing the air flow produced by the fan toward the PTC heater assemblies.\nAdditional objects, as well as features and advantages, of the present invention will be set forth in part in the description which follows, and in part will be obvious from the description or may be learned by practice of the invention. In the description, reference is made to the accompanying drawings which form a part thereof and in which is shown by way of illustration specific embodiments for practicing the invention. These embodiments will be described in sufficient detail to enable those skilled in the art to practice the invention, and it is to be understood that other embodiments may be utilized and that structural changes may be made without departing from the scope of the invention. The following detailed description is, therefore, not to be taken in a limiting sense, and the scope of the present invention is best defined by the appended claims."}
{"text": "Thermoelectric generators are becoming popular as backup sources of electricity to power electronic devices when traditional power sources are unavailable or wired grids go down. A thermoelectric generator includes a thermoelectric module with associated components for cooling the module, regulating power out, storing power and delivering power to various electronic devices. A variety of loads, such as heat-removing systems, energy storage devices, and electronics to be operated may be integrated to draw upon the electric energy created. Another aspect is availability of operational data such as thermoelectric module data and sensor data that could be harvested with a data logging system and stored in a data processing device and locally or remotely accessed. Thermoelectric generators work more efficiently when coupled to high thermal output devices such as wood and gas stoves, fireplaces, etc.\nTraditional thermoelectric generators rely on matching the output voltage of the thermoelectric generator to that of the load. The voltage and current of the generator come up together when connected to a heat source. Only after the output potential of the generator has achieved a potential greater than the resting voltage of the energy storage device or greater than the minimum voltage to power the load will it start to power the load. Most thermoelectric generators are designed to run on high temperature heat sources (i.e., 450° C. max temp). These generators struggle to even trickle charge an energy storage device at lower temperatures as their output voltage may be only a fraction of what would be necessary to charge the device. Having a system that can prioritize usage of generated electrical energy and convert overall electrical energy output from the thermoelectric module would be of benefit.\nThe present invention aims to provide a system for effectively charging energy storage devices at lower potential electrical energies while managing/prioritizing power used by all loads including active cooling devices for heat-exiting sides of the thermal electric generator. The system also regulates overloading of the thermoelectric modules and the active cooling system to prevent the components or parts of the internal and external system from being damaged."}
{"text": "Vehicle occupant protection systems having an actuatable protection element are well known in the art. One particular type of actuatable protection system includes an inflatable protection module. An inflatable protection device of the module is commonly referred to as an air bag, and is mounted to inflate within the occupant compartment of the vehicle. The protection module has a source of inflation fluid and an electrically actuatable igniter, referred to as a squib.\nThe occupant protection system further includes a crash sensor for sensing the occurrence of a vehicle crash condition and for providing an electrical signal indicative of the crash condition. When the sensor indicates that the vehicle is in a crash condition requiring inflation of the air bag to help protect a vehicle occupant (\"a deployment crash condition\"), an electric signal is provided to the squib for the purpose of igniting the squib. The squib, when ignited, activates the source of inflation fluid (e.g., ignition of a combustible gas or heat generating composition and/or opening of a container of pressurized gas). The source of inflation fluid is operatively coupled to the air bag, and, when activated, inflates the air bag.\nSeveral known occupant protection systems include an occupant position sensor and a controller, which controls the associated protection module in response to a sensed position of a vehicle occupant. In response to the sensed occupant position, the timing of air bag deployment, the pressure of the inflated air bag, aiming of the air bag, and/or disabling the air bag are controlled. These controllable/adjustable aspects, collectively, adjust the dynamic profile of the air bag. A protection system with controllable/adjustable aspects is commonly referred to as a \"smart\" system.\nOne example of a system with an adjustable dynamic profile is disclosed in U.S. Pat. No. 5,232,243 to Blackburn et al., and assigned to TRW Vehicle Safety Systems Inc. This patent discloses that it is not always desirable to inflate an air bag with 100 percent of the inflation fluid (e.g., gas) provided from an inflation fluid source. The system disclosed by Blackburn et al. controls the amount of gas that inflates the air bag in response to the detected weight of the occupant. Other example is disclosed in U.S. Pat. No. 5,330,226 to Gentry et al., and assigned to TRW Vehicle Safety Systems Inc. The system disclosed by Gentry et al. controls the amount of gas that inflates an air bag in response to detected occupant position."}
{"text": "This invention relates to converters, and more particularly to a nondissipative circuit for equalized load sharing of power transistors.\nThe present approach for securing reliability of operation throughout the life of a spacecraft mission for the electrical power systems is to use small converter modules operating in parallel to achieve the required system rating. As is pointed out in a copending application Ser. No. 684,809 (now abandoned), titled METHOD AND APPARATUS FOR AUTOMATIC LOAD SHARING AMONG PARALLELED CONVERTERS, by Col. W. T. McLyman and G. W. Wester, new systems can be tailored to new requirements by adding or removing modules without the need for developing new circuit designs and incurring the attendant development costs. However, as also pointed out in that application, parallel operation often results in very large differences in load sharing among the modules and unbalanced operation of the power transistors of each individual module. This occurs because of the low impedance characteristics of each converter module, and the significant effect of slight parameter differences in components, such as the power transistors, power transformers, etc. Such inequalities and unbalanced operation can cause either module or component level overstresses which may produce premature failure.\nThe aforesaid copending application describes a method for equalizing load sharing in duty-ratio operated converter modules by appropriately adjusting the duty ratio of each module as a function of the differential load currents provided by each. The scheme described therein takes advantage of the fact that since the converter modules are driven by pulse width modulators, adjusting the pulse width of each individual converter modulator changes the duty cycle to effect the desired sharing of the load by all the converter modules. However, since inverter operation must always be with a fixed duty cycle of 50%, that technique cannot be used in systems using inverter modules. Also that technique does not achieve balanced operation of the push-pull transistors of the power output stage. The invention disclosed herein overcomes those limitations in a circuit for equalizing load sharing of power transistors in a converter operated as a DC-AC inverter or a DC-DC converter.\nAs will become apparent from the following description, a converter embodying the present invention may be operated to convert direct current into alternating current, i.e., operated as an inverter, or to convert direct current into direct current. Thus, although the generic term \"converter\" has come to be used almost exclusively for DC-DC converters, and a special term \"inverter\" has come into general use for DC-AC converters, the generic term \"converter\" will be employed hereinafter to encompass both DC-DC and DC-AC converters. However, because the embodiment disclosed to illustrate the concept of the invention in a specific embodiment is a DC-AC converter, it will be referred to by the special term \"inverter.\""}
{"text": "1. Field of the Disclosure\nThe present invention relates in general to drill strings and, in particular, to a system, method and apparatus for regulating fluid flow through a drill string.\n2. Description of the Related Art\nConventional oil and gas drilling typically includes pumping a quantity of fluid through a pipe or drill string to a drill bit for cutting the hole in the rock. The fluid is then circulated back up though the wellbore in the annular or outer section of the hole. Drilling fluid is beneficial to the drilling process since it clears away pieces of rock that have been cut from the bottom of the wellbore. Without this cleaning action the cut pieces of rock would accumulate near the drill bit and interfere with further drilling.\nIn general, the higher level of fluid flow that a drilling operation can achieve, the better that cut pieces of rock or “cuttings” are cleared from the bottom of the wellbore. However, there are several factors that limit the fluid flow level. One of these factors is the amount of pressure that it takes to pump a large amount of fluid. As the drill string becomes longer or narrower, the resistance to pumping a given amount of fluid increases, which increases the need for higher pressure. With any fluid pump set up there is a limit to the amount of pressure that can be overcome in order to make the fluid flow. Accordingly, the size or type of pump can limit the available flow rate.\nAnother limiting factor is the capability of the downhole mud motor. Mud motors are used to make the rock cutting drill bit rotate faster than the drill pipe that it is connected to. For example, a drilling operator may desire to drill while holding the drill string stationary, or may want to rotate the drill bit faster to achieve a higher rate of rock penetration. The mud motor works in a manner similar to a turbine in that the mud that flows through the motor turns a rotor that is connected to the drill bit. Energy from the pressure of the fluid flow is converted into rotational work by the drill bit. Mud motors are usually designed such that there is a maximum amount of flow that the motors are designed to handle. Forcing excess fluid through a mud motor can damage the motor and inhibit the drilling process.\nThe desire to flow higher volumes of drilling fluid through the well and the need to limit the volume flow rate due to the constraints of the motor can be conflicting. It would be desirable to flow as much fluid as is desired while ensuring that the motor did not experience a rate of flow higher than its design criteria.\nA conventional solution to this problem is to form annular ports in the drill string above the mud motor. By choosing the size of the ports, the amount of flow that exits through the ports and the amount of flow that continues on through the drill string into the mud motor can be approximated.\nA problem with this technique is that the amount of fluid that exits through the ports varies depending on the back pressure from the mud motor. The back pressure from the mud motor is a factor of the torque that it delivers. Thus, the more torque that is needed or generated by the motor, the higher the back pressure from the motor, which diverts more fluid through the ports in the sides of the drill string. More diverted flow means less fluid is transferred down through the motor. Less fluid to the motor reduces its torque and power, which can induce a situation where the motor stalls and needs more torque to overcome its bound condition. Conversely, an off-bottom situation where there is relatively low amounts of back pressure generated by the motor because there is no drilling torque resistance can result in a higher amount of fluid passing through the motor and a lower amount of fluid exiting the drill string. This too is problematic since a low torque situation causes the motor to spin faster at a given flow rate. Increased amounts of flow will only exacerbate this situation.\nSome motor manufacturers attempt to solve this problem by drilling a hole through the rotor of the mud motor so that some fluid may pass through the tool without generating torque or causing damage to the motor. Unfortunately, since the drilled hole is static and does not change its shape to account for differing flow or pressure conditions, it is subject to the same limitations as the previously described method. Thus, improvements in controlling drill string fluid flow continue to be of interest."}
{"text": "The present invention relates to a router bit. A router bit used for chamfering and finishing the butt end and side of a board is available in two types, i.e. a straight tooth type and a helical tooth type. Usually, a router bit of the helical tooth type is provided with an even number of cutting edges having helix angles in the same hand or direction with respect to the axis of rotation of the cutting edges. The conventional router bit of this type is liable to form burrs at the side of board when a board material containing fiber reinforced plastics is chamfered or finished. This is true particularly in chamfering or finishing FRP board or FRP-clad board in which fluffy fibers remain on the finished surface because the fibers are much stronger than the matrix resin. Such rough finish brings about delamination, corrosion, a decrease of strength, and breakage. This is quite different from problems involved in chamfering or finishing metal materials.\nThe object of the present invention is to eliminate the drawbacks of conventional technology as mentioned above, and is the result of observation of mutual actions of cutting edges and the material being cut."}
{"text": "The present invention relates to a toner and a process for producing the same that is used for developing an electrostatic latent image formed by an electrophotographic method and an electrostatic recording method, a developer for an electrostatic latent image, and a process for producing an image.\nA method of visualizing image information via an electrostatic image, such as an electrophotographic method, is utilized in various field of art. In the electrophotographic method, an electrostatic static latent image is formed on a photoreceptor through charging and exposure steps, and the electrostatic latent image is developed with a developer containing a toner, which is then visualized through transfer and fixing steps.\nAs the developer used herein, a two-component developer containing a toner and a carrier and a one-component developer solely using a magnetic toner or a non-magnetic toner are known.\nThe production of the toners are generally conducted by a kneading and pulverization method, in which a thermoplastic resin is melted and kneaded with a pigment, a charge controlling agent and a releasing agent, such as a wax, and after cooling, it is finely pulverized and classified. In order to improve the fluidity and the cleaning property, inorganic or organic fine particles are sometimes added to the surface of the toner particles depending on necessity. While a toner of good quality can be produced by the process described, the following problems are involved.\nThe shape of the toner and the surface structure of the toner do not become constant by using the ordinary kneading and pulverization method, and they are delicately changed depending on the pulverization property of the material used and the pulverization conditions. Therefore, it is difficult to positively control the shape of the toner and the surface structure of the toner. Furthermore, there is restriction on the selection of the material in the kneading and pulverization method. Specifically, it is necessary that the dispersion body of the resin containing the coloring agent dispersed therein is sufficiently brittle and can be finely pulverized by an apparatus that can be practically employed from the standpoint of economy.\nHowever, when the dispersion body of the resin and the coloring agent is brittle, there are cases where fine powder is formed or the shape of the toner is changed due to a mechanical sharing force in a developing device. As a result, the fine powder is fixed on the surface of the carrier to accelerate deterioration of the charging property of the developer in the two-component developer, and the particle size distribution is broadened to cause the toner-smoke and the developing property is lowered by the change of the toner shape to deteriorate the image quality in the one-component toner.\nIn the case of toner to which a large amount of a releasing agent such as wax is internally added, the releasing agent is often exposed on the surface of the toner depending on the combination with the thermoplastic resin, so as to influence the property of the toner. Particularly, in the case of a combination of a resin that is relatively difficult to be pulverized owing to a high elasticity due to a high molecular weight component and brittle wax such as polyethylene, exposure of the polyethylene on the toner surface is often observed. This phenomenon is advantageous for the releasing property on fixing and the cleaning of a non-transferred toner on the photoreceptor, but the polyethylene on the surface layer is easily transferred by a mechanical force to contaminate the developing roll, the photoreceptor and the carrier, which brings about reduction in reliability.\nFurthermore, when the shape of the toner is irregular, the fluidity cannot be sufficiently ensured even though a fluidizing aid is added, and thus the fine particles on the surface of the toner migrate to concave parts of the toner by a mechanical sharing force upon use to lower the fluidity with the lapse of time, or the fluidizing aid is buried inside the toner to deteriorate the developing property, the transfer property and the cleaning property. Furthermore, in the case where a toner recovered by cleaning is reused by bringing back to the developing device, the image quality is liable to be lowered. When the amount of the fluidizing aid is increased to prevent the problems, other problems occur in that black spots are formed on the photoreceptor, and the particles of the aid are scattered.\nIn recent years, as a method for positively controlling the shape and the surface structure of the toner, a process for producing a toner by an emulsion polymerization and aggregation method has been proposed in JP-A-282752/1988 and JP-A-250439/1994. In the method, a resin fine particle dispersion is firstly prepared by emulsion polymerization, and separately, a coloring agent dispersion is prepared by dispersing a coloring agent in a solvent, followed by mixing and heating both the dispersion to produce a toner through fusing and coalescence. While the shape of the toner can be controlled in a certain extent to improve the charging property and the durability in this process, because the internal structure becomes substantially uniform, there are problems in the releasing property of a fixing sheet on fixing and in the stability of transparency on outputting on an OHP sheet.\nAs described in the foregoing, in the electrophotographic process, in order to stably maintain the performance of the toner even under various kinds of mechanical stress, it is necessary that the exposure of the releasing agent on the toner surface is controlled, and the surface hardness is increased without deteriorating the fixing property, so as to increase the mechanical strength of the toner itself and to sufficiently maintain both the charging property and the fixing property.\nOwing to increase in demand for high image quality in recent years, there are considerable trends of reduction in particle diameter of the toner to realize an image of high definition particularly in the formation of a color image. However, when the particle diameter is simply decreased maintaining the conventional particle size distribution, problems of contamination of the carrier and the photoreceptor and the toner-smoke due to the presence of the toner as fine particles become outstanding, and it is difficult to simultaneously realize both the high image quality and the high reliability. In order to solve the problem, it becomes important that the particle diameter is decreased, and simultaneously the particle size distribution is sharpened.\nIn digital full-color duplicator and printer in recent years, a color image is subjected to color separation by filters of B (blue), R (red) and G (green), and a latent image formed with dots having a diameter of from 20 to 70 xcexcm corresponding to the original image is developed by using toners of Y (yellow), M (magenta), C (cyan) and Bk (black) with utilization of the subtractive color mixing. In comparison to an ordinary monochrome duplicator, it is necessary that a full-color duplicator transfers a large amount of a toner and handles a small dot diameter, and therefore the importance of the uniform charging property, the stability, the toner strength and the sharpness of the particle size distribution is increased. In order to meet the high-speed performance and the energy saving property of the apparatus, the further low temperature fixing property is demanded. Therefore, the aggregation and fusing/coalescence process which is suitable to produce such toner with sharpness of particle size distribution and with small particle diameter is preferred to impart excellent characteristics to the toner.\nBecause a large amount of toners is installed in a full color duplicator, it is necessary to sufficiently conduct color mixing thereof, and improvement in color reproducibility and transparency on OHP on color mixing becomes necessary.\nAs a component of the releasing agent, polyolefin series wax is generally used to prevent low temperature offset upon fixing, which is internally added to the toner. In addition to this, a slight amount of a silicone oil is uniformly coated on a fixing roller to improve the high temperature offset property. Therefore, there is a problem of causing uncomfortable sticky feeling due to attachment of the silicone oil to the transfer material.\nJP-A-61239/1993 proposes a toner suitable for oilless fixing, in which no oil is coated on a fixing roll, by incorporating a large amount of a releasing agent component in the toner. However, when a large amount of releasing agent is added, the releasing property can be certainly improved, but the binder resin component and the releasing agent are dissolved in each other, and oozing of the releasing agent becomes non-uniform, so as to damage the stability of releasing. Furthermore, there are cases where a free component of the releasing agent causes charging inhibition.\nAs methods for solving the problems, a method of adding a high molecular weight component to increase the rigidity of the binder resin is proposed in JP-A-69666/1992 and JP-A-258481/1997, and a method of introducing crosslinking to the binder resin to supplement the rigidity is proposed in JP-A-218460/1984 and JP-A-218459/1984.\nHowever, when a crosslinking component is added to the binder resin in the later case, because the rigidity of the binder resin is increased, the flexibility of a fixed image becomes poor. Particularly, when only the molecular weight of the crosslinking agent is increased in JP-A-218460/1984, the molecular weight between the crosslinking points is increased to slightly improve the flexibility of the fixed image itself, but because both the molecular weight and the crosslinking density of the binder resin are not simultaneously controlled, it is difficult that the releasing property, the gloss of the surface of the fixed image and the transparency on OHP cannot be simultaneously attained with maintaining the sufficient flexibility. In particular, a sufficient fixed image cannot be obtained when it is used in a fixing device of an energy saving type or a duplicator or a printer of a high printing speed disclosed in JP-A-221987/1998.\nThe invention has been made to solve the problems and to provide a toner for developing an electrostatic latent image, a process for producing the same, a developer for an electrostatic latent image, and a process for forming an image in that the adhesion property of a fixed image to a fixing sheet, the peeling property of a fixing sheet, the hot offset resisting property, the flexibility of a fixed image, the gloss of the surface of the fixed image and the transparency on OHP are excellent, the uniformity and the stability in charging are high, fogging and scattering of the toner are caused, and an image of high quality can be formed.\nThe inventors have been conducted earnest investigations mainly in the molecular weight between crosslinking points Mc and the crosslinking density Me, and as a result, it has been found that the problems are solved by the following configurations.\nThe present invention provides a toner for developing an electrostatic latent image containing a binder resin and a coloring agent, the toner having a molecular weight between crosslinking points Mc obtained by a measurement of temperature dispersion on dynamic viscoelasticity of from 3.6xc3x97106 to 7.5xc3x97108, and the molecular weight between crosslinking points Mc and a crosslinking density Me satisfying the relationship: 14.0xe2x89xa6Log10(Mc/Me)xe2x89xa616.5. The toner may have a glass transition temperature Tg of from 45 to 52xc2x0 C. The toner may have a weight average molecular weight in a range of from 25,000 to 55,000. The toner may have a volume average particle size distribution index GSDv of 1.30 or less, and a ratio of the volume average particle size distribution index GSDv and a number average particle size distribution index GSDp of 0.95 or more. The toner may have a shape factor SF1 in a range of from 100 to 140. The coloring agent contained in the toner may have a central diameter in a range of from 100 to 330 nm measured by a transmission electron microscope (TEM). An amount of the coloring agent in the toner may be in a range of from 4 to 15% by weight based on a total weight of the toner. The toner may further contains a releasing agent. An amount of the releasing agent in the toner may be in a range of from 5 to 25% by weight based on a total weight of the toner. The toner may have an absolute value of a charging amount of from 20 to 40 xcexcC/g.\nThe present invention also provides a process for producing a toner for developing an electrostatic latent image, the process containing a step of mixing a resin particle dispersion containing resin particles dispersed therein and a coloring agent dispersion containing a coloring agent dispersed therein to aggregate fine particles, and a step of fusing the aggregated particles by heating. The process may further contains, after the aggregationn step, a step of additionally mixing the resin particle dispersion. The step of aggregating the fine particles may contain mixing a releasing agent dispersion. A polymer of a metallic salt may be used in the step of aggregating the fine particles.\nThe toner may be used as a developer with a carrier particles.\nThe present invention also provides a process for producing an image containing a step of developing an electrostatic latent image formed on a latent image holding member with the developer, a step of transferring a toner image thus formed to a transfer material, and a step of fixing the toner image to the transfer material. The process may further contains, after transferring step, a step of recovering the toner remaining on the latent image holding member, and a recycle step of bringing the recovered toner to the developer. A cleaning step may not be conducted after the transferring step. The fixing step may be conducted by using a belt in the fixing step."}
{"text": "The present invention relates generally to the production of relativistic particles, and, more particularly, to the production of relativistic particles by colliding a plasma with electromagnetic pulses.\nThe sustainable acceleration of charged particles is a major challenge in both basic plasma physics and astrophysics. Current conventional particle accelerators have an energy gain per distance of only about 1GeV/m. Recent advances in ultra-intense lasers (ULs) (e.g., with intensity I>2×1018Wcm−2) open up a new frontier in the acceleration of particles via intense electromagnetic (EM) fields. (See e.g., G. A. Mourou, C. P. J. Barty, M. D. Perry, Phys. Today 51(1), 22 (1998), the SAUUL Report, T. Ditmire, Ed. (UT Austin, 2003), for and by L Lontano et al., Eds., Superstrong Fields in Plasmas, AIP Conf. Proc. No. 611 (AIP, NY 2002)).\nMost proposed laser acceleration schemes, for example, Laser Wake Field Acceleration (LWFA), Plasma Wake Field Acceleration (PWFA), Plasma Beat Wave Acceleration (PBWA), as described by E. Esarey, P. Sprangle, J. Krall, A. Ting, IEEE Trans. Plasma Sci. 24, 252 (1996), P. Sprangle, E. Esary, A. Ting, Phys. Rev. Lett. 64, 2011 (1990), V. Malka, in AIP Conf. Proc. No. 611, p. 303, Ed. M. Lontano et al. (AIP, NY, 2002), A. Pukhov, J. Meyer-ter-Vehn, Phys. Rev. Lett. 79, 2686 (1997), T. Tajima and J. M. Dawson, Phys. Rev. Lett. 43, 267 (1979), and Free Wave Acceleration (FWA), as described by M. S. Hussein, M. P. Pato, A. K. Kerman, Phys. Rev. A 46, 3562 (1992), M. S. Hussein, M. P. Pato, Phys. Rev. Lett. 68, 1992, S. Kawata, T. Maruyama, H. Watanabe, I. Takahashi, Phys. Rev. Lett. 66, 2072 (1991), J. G. Woodworth, M. N. Kreisler, A. K. Kerman, The Future of Accelerator Phys. p. 378, Ed. T. Tajima (AIP, N.Y. 1996), involve the propagation of lasers in an underdense plasma (ωpe=(4πne2/me) 1/2<ωo=2πc/λ, where λ=laser wavelength, me is the rest mass of an electron e− or a positron e+, e is the absolute value of the electric charge of an electron e− or a positron e+, c is the speed of light in a vacuum, and n=electron density). In such schemes both the acceleration gradient (energy gain/distance) and the particle beam intensity are limited by the underdense plasma requirement. (See e.g., E. Esarey, P. Sprangle, J. Krall, A. Ting, IEEE Trans. Plasma Sci. 24, 252 (1996), P. Sprangle, E. Esary, A. Ting, Phys. Rev. Lett. 64, 2011 (1990), V. Malka, in AIP Conf. Proc. No. 611, p. 303, Ed. M. Lontano et al. (AIP, NY, 2002), A. Pukhov, J. Meyer-ter-Vehn, Phys. Rev. Lett. 79, 2686 (1997), T. Tajima and J. M. Dawson, Phys. Rev. Lett. 43, 267 (1979)).\nWhen a single ultra-intense laser (UL) irradiates a plasma surface through a variety of nonlinear collisionless processes, as described, for example, by W. L. Kruer, E. J. Valeo, K. G. Estabrook, Phys. Rev. Lett. 35, 1076 (1975), W. L. Kruer, K. G. Estabrook, Phys. Fluids 28, 430 (1985), and S. C. Wilks, W. L. Kruer, M. Tabak, A. B. Langdon, Phys. Rev. Lett. 69, 1383 (1992), the single ultra-intense laser (UL) as described, for example, by S. C. Wilks, W. L. Kruer, M. Tabak, A. B. Langdon, Phys. Rev. Lett. 69, 1383 (1992), and A. Pukhov, J. Meyer-ter-Vehn, Phys. Rev. Lett. 79, 2686 (1997), couples a significant fraction of the energy of the single ultra-intense laser (UL) to superthermal electrons with characteristic energy E given by the Lorentz relativistic gamma factor γ=E/(511 keV)˜(1+Iλ2/1.4×1018Wcm−2)1/2−1, where I is the laser intensity in Watts/cm2 and λ is the laser wavelength in μm,.\nIf the plasma is a slab of electron-positron (e+e−) pairs, in addition to the collisionless heating, the light pressure will also snowplow the e+e− pairs with a bulk velocity determined by momentum conservation, as described, for example, by W. L. Kruer, E. J. Valeo, K. G. Estabrook, Phys. Rev. Lett. 35, 1076 (1975). For λ=1 μm and I=1021W/cm2, particle-in-cell (PIC) simulations, which are described, for example, by S. C. Wilks, W. L. Kruer, M. Tabak, A. B. Langdon, Phys. Rev. Lett. 69, 1383 (1992), A. Pukhov, J. Meyer-ten-Vehn, Phys. Rev. Lett. 79, 2686 (1997), and E. P. Liang, S. C. Wilks, M. Tabak, Phys. Rev. Lett. 81, 4887 (1998), show that the electrons can be accelerated to greater than about 10MeV. This has been confirmed with experiments, as described, for example, by S. P. Hatchett et al. Phys. Plasmas 7, 2076 (2000), K. W. D. Ledingham et al. Phys. Rev. Lett. 84, 899 (2000). and T. E. Cowan et al., Phys. Rev. Lett. 84, 903 (2000).\nIn a conventional laser ponderomotive accelerator in which a single ultra-intense laser (UL) strikes an overdense e+e− plasma surface, as described, for example, by S. C. Wilks, W. L. Kruer, M. Tabak, A. B. Langdon, Phys. Rev. Lett. 69, 1383 (1992), all upstream particles share the momentum of the Poynting flux. For a laser wavelength λ=1 μm and an ultra-intense laser (UL) intensity I=1021W/cm2, in conventional laser ponderomotive heating, the Lorentz relativistic gamma factor γ may reach only γ˜Ωe/ωo˜30, where Ωe=eB/mec is equal to the electron gyrofrequency in the laser magnetic field B, me is the rest mass of an electron e− or a positron e+, e is the absolute value of the electric charge of an electron e− or a positron e+, c is the speed of light in a vacuum, and ωo is equal to the laser frequency. (See e.g., W. L. Kruer, K. G. Estabrook, Phys. Fluids 28, 430 (1985) and S. C. Wilks, W. L. Kruer, M. Tabak, A. B. Langdon, Phys. Rev. Lett. 69, 1383 (1992)). A Diamagnetic Relativistic Pulse Accelerator (DRPA) of an overdense plasma has been proposed, as described, for example, by E. Liang, K. Nishimura, H. Li, S. P. Gary, Phys. Rev. Lett. 90, 085001 (2003), E. Liang, K. Nishimura, Phys. Rev. Lett. 92, 175005 (2004), and K. Nishimura, E. Liang, Phys. Plasmas, 11 (10) (2004). However, the Diamagnetic Relativistic Pulse Accelerator (DRPA) is difficult and expensive to realize practically.\nHigh-energy gamma-ray (γ-ray) beams are conventionally only produced by large conventional particle accelerators (e.g., synchrotron sources) at national research facilities. Both the intensity of the high-energy gamma-ray (γ-ray) beams produced by the large conventional particle accelerators, and the energy conversion efficiency, are relatively low. Consideration is now being given to the design of particle accelerators that can produce relativistic particles. In particular, attention is being directed to the production of relativistic particles by colliding a plasma with electromagnetic pulses."}
{"text": "The present invention relates to an electron beam probe system which irradiates an IC under test with an electron beam, measures the quantity of secondary electrons emanating from each irradiated point and displays the potential distribution in the IC as a potential contrast image that is used to locate a fault, for example. The invention also pertains to an IC analysis system and a fault isolation method that utilize such an electron beam probe system."}
{"text": "This invention relates generally to the field of cat toys and more specifically to an impact action cat toy.\nCat's who live in a home environment enjoy practicing their hunting skills by chasing after erratically moving items whether they be real—such as mice, or artificial, such as a cat toy that is operated by the cat's owner to imitate the erratic action of a mouse or other naturally occurring small animal.Numerous cat toys have been designed and manufactured to fill this need.However, most of those toys do not provide a high degree of instantaneous and quick movement that cats love so much. Simple cat toys that dangle from a string can be dragged and flown by the user but require a great deal of effort on the user's part, and even then, the action is not necessarily lightning fast because a person's hand and arm can only move to a certain degree of speed and no further. What is needed is a cat toy that includes a novel way to produce a truly instantaneous movement and that is very simple and economical to manufacture."}
{"text": "The present invention relates to a transmission control method in a radio communication system. Particularly, the represent invention relates to a transmission control method in a communication system that establishes packet communications where a radio communication medium is shared between plural terminals and a base station and that uses a multiple access system in which possible collision occurs.\nConventionally, in packet communications, plural terminals often share the communication medium because of burst due to packet occurrence. In such a case, the shared communication medium provides a point-to-multipoint transmission line. If the communication medium is wireless, the forward link which transmits signals from a radio base station to terminals is used for broadcasting so that all terminals become a receivable state. The reverse link which transmits signals from a terminal to the radio station is generally used as a multi-access in which plural terminals use a common communication medium. To provide an effective multi-access mode to the reverse link means to increase the transmission capacity per shared communication medium. As a result, this allows the subscriber capacity of the system to be greatly increased.\nTo solve such a problem, various multi-access systems have been conventionally proposed. The basic system which can provide the simplest control is the ALOHA system (refer to Telecommunication Networks (ISBNO-201-16423), pp. 407-408). According to the basic concept of the system, packet signals are transmitted immediately when they are newly generated. This allows a very small delay to be achieved although collision may occur. This access system in which possible collision occurs is hereinafter called a collision-based access system.\nMoreover, there is the CSMA (Carrier Sense Multiple Access) system (refer to Telecommunication Networks (ISBNO-201-16423), pp. 437) which detects whether or not other terminal currently is in a transmission mode by means of a receiving section and then transmits packet signals immediately when the idle period detection. This system allows collision of transmitted packets to be relieved because packet signals are not transmitted while other terminal is in a transmission mode.\nHowever, when the case where a small number of terminals in mobile communications exits in a line-of-sight path and the effect of a hidden terminal is large, the terminal cannot completely sense the carrier. In such a case, the ICMA system (Idle-signal Casting Multiple Access) system is used to reduce collision of a packet signal at a multi-access time, thus improving the throughput.\nMoreover, in consideration of survival at a packet signal collision time due to a capture effect, the ICMA-PE (Idle-signal Casting Multiple Access With Partial Echo) system that uses a portion of an up signal, or a partial echo, as an access control signal (Umeda and Onoe, xe2x80x9cIdle-signal Casting Multiple Access With Partial Echoxe2x80x9d Communication Technical Reports RCS91-30) has been proposed. FIG. 4 is a sequence diagram showing the ICMA-PE system.\nReferring to FIG. 4, when the radio base station 40 receives an reverse link message Dlu(0) from the terminal 41, it sets up a busy flag as access control information and adds a portion of the received data Dlu(0) as a partial echo PE1(0). This operation can determine whether or not which terminal has transmitted a survival signal due to packet signal collision by a capture effect, thus enabling a relief process. Moreover, as a simple ICMA-PE system, a multiple access system called the ICMA-CD (Idle-signal Casting Multiple Access With Collision Detection) has been proposed.\nHowever, in the collision-based access system, since collision cannot be avoided at the beginning of transmission, a serious problem occurs in the case where many terminals want transmission. Moreover, since the radio station cannot judge that a certain transmission error has occurred due to collision or degradation in the quality of a radio line, it is difficult to deal with such a problem.\nThere is an access system in which collision does not occur different from the above-mentioned collision-based access system. The multiple access system is hereinafter called the non-collision-based access system. This access system is generally realized by circulating the transmission right of a terminal. There is a multiple access system called polling system as one of non-collision-based access systems.\nIn the polling system, a radio base station asks each terminal whether there is a signal to be transmitted. If there is a signal to be transmitted to a terminal, the terminal transmits a signal so that polling (transmission invitation) is sequentially performed. In this polling system, since the radio station can perform centralized management, the signal collision at a multiple access time do not occur between plural terminals. However, even if a signal to be transmitted occurs at a terminal, the signal cannot be transmitted until polling invests it with a transmission right. Hence, there is the problem that the transmission delay occurs.\nAs one approach to solve such a problem, JP-A-274788/1996 discloses a multiple access method that uses adaptively and selectively the collision-based access system and the non-collision-based access system. That is, according to the communication traffic between a terminal and a radio base station, the connection delay is decreased by using the collision multiple access system at the time of low traffic and using the non-collision multiple access system at the time of high traffic. Moreover, high throughput can be realized even at a high traffic time.\nAccording to the invention disclosed in JP-A-274788/1996, the collision-based access system employs such a methodology in which the terminal that judges that a packet signal transmitted from itself to avoid re-collision after the occurrence of collision has collided is transferred to a back-off state and then re-transmits the packet signal after a period of a random time or transmits the next signal. On the other hand, in the non-collision-based access system, since a terminal can perform transmission only by the polling from a radio base station, collision does not occur, so that the terminal does not change to a back-off state.\nAs described above, in the multiple access method adaptively using the collision-based access system and the non-collision-based access system, when the mode changes from the collision-based access system to the non-collision-based access system due to collision of a signal transmitted from a certain terminal, the terminal is changed to a back-off state. In this case, the problem arises that the terminal cannot perform transmission operation because of its back-off state even if polling is performed from the radio base station to the terminal.\nThe objective of the present invention is to solve the above-mentioned tasks.\nMoreover, the objective of the invention is to provide a transmission control method that can improve a decrease in throughput due to a transmission disabled state when a terminal is changed to a back-off state due to collision happened during the use of a collision access system in a communication system using at least a collision-system multiple access system, thus becoming a transmission wait state.\nIn order to accomplish the above mentioned objects, the objective of the present invention is achieved by the transmission control method suitable in a radio communication system which performs packet communications by sharing a radio communication medium between plural terminals and a radio station and employs a multiple access system in which collision may occur, the method comprising the steps of switching, a terminal to a transmission mode, or if the terminal is in transmission mode keeping it in transmission mode, during use of a collision-based access system when receiving a polling from said radio base station, said terminal being changed to a back-off state, thus being in a wait state, because of collision of a signal transmitted from said terminal.\nIn the transmission control method according to the present invention, each of the terminals judges that a packet signal transmitted from its terminal has collided when its terminal receives access control information from the radio base station after transmitting the, packet signal to the radio base station, the access control information representing no reception, and changes to a back-off state in which a different random time transmission is suspended.\nIn the transmission control method according to the present invention, the radio base station transmits a common polling signal to all the terminals when a packet signal from each of the terminals is not received, and sequentially transmits polling signals to other terminals after a packet is completely received from a specific terminal.\nFurthermore, in the transmission control method suitable in a radio communication system which performs packet communications by sharing a radio communication medium between plural terminals and a radio station and adaptively employs a multiple access system in which collision may occur and a multiple access system in which collision does not occur, the method comprises the steps of detecting collision by means of the radio station; changing the radio communication system to a back-off state by transmitting the collision signal when the multiple access system in which collision may occur is changed to the multiple access system in which collision does not occur; and making the terminal in a transmission mode at a polling receive time when the terminal in a wait state has received polling from the base station.\nIn the transmission control method according to the present invention, each of the terminals judges that a packet signal transmitted from itself has collided when receiving the access control information from the radio base station, the access control information representing that the packet signal is not received, after transmission of the packet signal to the radio station, and changes to a back-off state in which a different random time transmission is sustained.\nMoreover, in the transmission control method suitable in a radio communication system which performs packet communications by sharing a radio communication medium between plural terminals and a radio base station and employs a multiple access system in which collision may occur in a reverse link from the plural terminals to the radio base station and which realizes transmission confirmation regarding a forward link packet signal to be transmitted on a forward link in response to a response signal from a terminal being a destination of the forward link packet signal, the method comprises the steps of transmitting a collision signal after the radio base station detects collision; and transmitting a response signal for transmission confirmation at a polling receiving time when a terminal changed to a back-off state and in a wait state receives polling from the radio station.\nAccording to the present invention, when the radio communication system is changed from the collision-based access system to the non-collision-based access system or a response signal is sent back to a signal transmitted on the forward link, a terminal in a back-off state or wait state can be quickly changed to a transmission mode in response to a polling signal from a radio base station. As a result, the throughput is improved and the radio line can be effectively used. Moreover, a response signal to be transmitted from a terminal for transmission confirmation can be quickly invited to a signal to be transmitted from a base station to a terminal.\nAccording to the present invention, the transmission control method comprises the steps of transmitting a first forward link message from a radio base station to terminals; transmitting packet signals from the terminals in response to the first forward link message; transmitting a second forward link message from the radio base station when the radio base station does not receive the packet signals due to collision of the packet signals; changing the terminals to a back-off state in response to the second forward link message to avoid re-collision; retransmitting a packet signal from a terminal first recovered from the back-off state in response to a third forward link message from the radio base station; transmitting a fourth forward link message containing access control information from the radio base station first recovered terminal when the radio base station correctly receives the packet signal from the first recovered terminal; and sequentially transmitting a fifth forward link message from the radio base station when a packet signal from the first recovered terminal is transmitted.\nMoreover, according to the present invention, the transmission control method comprises the steps of transmitting a first forward link message from a radio base station to terminals; transmitting packet signals from the terminals in response to the first forward link message when the radio base station is in a collision-system access mode; changing the radio base station from a collision-system access mode to a non-collision-system access mode when the radio base station does not receive the packet signals due to collision of the packet signals, and then transmitting a second forward link message from the radio base station; changing the terminals to a back-off state in response to the second forward link message to avoid re-collision; retransmitting a packet signal from a terminal first recovered from the back-off state according to transmission invitation in response to a third forward link message from the radio base station; transmitting a fourth forward link message containing polling information from the radio base station to the first recovered terminal when the radio base station correctly receives the packet signal from the first recovered terminal; and sequentially transmitting a fifth forward link message from the radio base station when a packet signal from the first recovered terminal is transmitted.\nAccording to the present invention, the transmission control method comprises the steps of transmitting packet signals from terminals in response to a first forward link message from a radio base station; transmitting a second forward link message from the radio base station when the radio base station does not receive the packet signals due to collision of the packet signals; changing the terminals to a back-off state in response to the second forward link message to avoid re-collision; extracting forward link message for a terminal in a back-off state from a third forward link message transmitted from the radio base station retransmitting a packet signal from a terminal first recovered from the back-off state in response to the third forward link message from the radio base station; transmitting a fourth forward link message containing access control information and forward link message for first-recovered terminal from the radio base station when the radio base station correctly receives the packet signal from the first recovered terminal; extracting forward link message from the fourth forward link message; transmitting a fifth forward link message containing access control information and forward link message from the radio base station to confirm transmission of forward link message; and transmitting a response packet signal from the terminal in a back-off state in response to the fifth forward link message."}
{"text": "The Third Generation Partnership Project (3GPP) is standardizing a first release of a Long Term Evolution (LTE) concept. An LTE system or network may include a number of base stations (also referred to as “Node Bs”) connected to one or more mobility management entities/serving gateways (MME/S-GWs). A number of nodes (e.g., a network management system (NMS) node, an operation and support system (OSS) node, etc.) may provide operation and maintenance functionality for the base stations and/or the MME/S-GWs. In LTE, a downlink is based on orthogonal frequency division multiplexing (OFDM), while an uplink is based on a single carrier modulation method known as discrete Fourier transform spread OFDM (DFT-S-OFDM). During initial access, user equipment (UE) (e.g., a wireless telephone, personal digital assistant (PDA), etc.) seeks access to a network (e.g., a radio network) in order to register and commence services. A LTE random access (RA) procedure serves as an uplink control procedure to enable the UE to access the LTE network.\nThe LTE network may provide access to a variety of other networks, such as the Internet, Internet Protocol (IP) Multimedia Subsystem (IMS) networks, the Public Switched Telephone Network (PSTN), etc. In the LTE standard, when the user equipment is powered on it automatically attaches to the LTE network (i.e., the base stations, MME/S-GWs, etc.) even if there is no activity (e.g., a request for a service, a communication, etc.) from a user of the user equipment. Such unnecessary LTE network connectivity consumes and wastes resources associated with the LTE network."}
{"text": "Cellular radiotelephone system base transceiver stations (BTSs), at least for some United States (U.S.) and European Union (EU) applications, may be constrained to a maximum allowable effective isotropically radiated power (EIRP) of 1640 watts. EIRP, as a measure of system performance, is a function at least of transmitter power and antenna gain. As a consequence of restrictions on cellular BTS EIRP, U.S., EU, and other cellular system designers employ large numbers of BTSs in order to provide adequate quality of service to their customers. Further limitations on cells include the number of customers to be served within a cell, which can make cell size a function of population density.\nOne known antenna installation has an antenna gain of 17.5 dBi, a feeder line loss of 3 dB (1.25″ line, 200 ft mast) and a BTS noise factor of 3.5 dB, such that the Ga−NFsys=17.5−3.5−3.0=11 dBi (in uplink). Downlink transmitter power is typically 50 W. With feeder lines, duplex filter and jumper cables totaling −3.5 dB, the Pa input power to antenna is typically 16 W, such that the EIRP is 16 W+17.5 dB=1,000 W.\nIn many implementations, each BTS is disposed near the center of a cell, variously referred to in the art by terms such as macrocell, in view of the use of still smaller cells (microcells, nanocells, picocells, etc.) for specialized purposes such as in-building or in-aircraft services. Typical cells, such as those for city population density, have radii of less than 3 miles (5 kilometers). In addition to EIRP constraints, BTS antenna tower height is typically governed by various local or regional zoning restrictions. Consequently, cellular communication providers in many parts of the world implement very similar systems.\nRestrictions on cellular BTS EIRP and antenna tower height vary within each countries. Not only is the global demand for mobile cellular communications growing at a fast pace, but there are literally billions of people, in technologically-developing countries such as India, China, etc., that currently do not have access to cellular services despite their willingness and ability to pay for good and inexpensive service. In some countries, government subsidies are currently facilitating buildout, but minimization of the cost and time for such subsidized buildout is nonetheless desirable. In these situations, the problem that has yet to be solved by conventional cellular network operators is how to decrease capital costs associated with cellular infrastructure deployment, while at the same time lowering operational expenses, particularly for regions with low income levels and/or low population densities. An innovative solution which significantly reduces the number of conventional BTS site-equivalents, while reducing operating expenses, is needed."}
{"text": "The (Al,Ga,In)N material system includes materials having the general formula AlxGayIn1-x-yN where 0≦x≦1 and 0≦y≦1. In this application, a member of the (Al,Ga,In)N material system that has non-zero mole fractions of aluminium, gallium and indium will be referred to as AlGaInN, a member that has a zero aluminium mole fraction but that has non-zero mole fractions of gallium and indium will be referred to as InGaN, a member that has a zero indium mole fraction but that has non-zero mole fractions of gallium and aluminium will be referred to as AlGaN, and so on. There is currently considerable interest in fabricating semiconductor light-emitting devices in the (Al,Ga,In)N material system since devices fabricated in this system can emit light in the blue-violet wavelength range of the spectrum (corresponding to wavelengths in the range of approximately 380-450 nm).\nFIG. 1 is a schematic sectional view of a semiconductor laser device. The laser device consists of a laser structure 2 that is grown over a substrate 1. In this example, the laser device is a separate confinement (SCH) laser device and the laser structure 2 comprises, in sequence, a first cladding region 4, a first optical guiding region 5, an active region for laser oscillation 6, a second optical guiding region 7, and a second cladding region 8. A buffer layer 3 is provided between the substrate 1 and the laser structure 2. The first cladding region 4 and the first optical guiding region 5 are doped n-type, and the second cladding region 8 and the first optical guiding region 7 are doped p-type. A p-type contact layer 9 is disposed on the p-type cladding region 8. The laser further has a second, n-type contact, and this may be provided on the same surface of the substrate 1 as the laser structure (as indicated at 10a) or it may be provided on the opposite surface of the substrate 1 to the laser structure (as indicated at 10b).\nThe first cladding region 4, first optical guiding region 5, active region 6, second optical guiding region 7, second cladding region 8 and the contact layer 9 have been etched to a desired width w, to form a mesa. The mesa may be a ridge mesa and extend into the plane of the paper, or the mesa may have circular symmetry about the vertical axis. The width w of the mesa will be small, typically a few μm, so that the surface area of the mesa is small (particularly in the case of a circularly symmetric mesa). In order to ensure good electrical contact to the laser it is therefore conventional to deposit an electrically conductive “pad” 11 over the laser. The conductive pad 11 has a much greater surface area than the mesa, and so it is easier to make an electrical contact to the pad 11 than it is to make an electrical contact direct to the contact layer 9.\nProviding the electrically conductive pad 11 directly over the laser structure, as shown in FIG. 1, has the disadvantage that the conductive pad 11 is disposed directly on the portions of the buffer layer 3 that are exposed in the step of etching to form the mesa, and so makes electrical contact with the buffer layer 3. As a result, low resistance current paths 12,13 exist that go from the electrically conductive pad 11 to the n-type contact 10a or 10b without passing through the active region 6. The current paths 12, 13 that do not pass through the active region 6 are in parallel to the desired current path 14 through the active region 6, and degrade the characteristics of the laser-current that flows along these current paths 12,13 does not contribute to laser oscillation. It is therefore usual to provide a current blocking layer 15 having a high electrical resistance between the buffer layer 3 and the conductive pad 11, as shown in FIG. 2. The laser of FIG. 2 is generally similar to the laser of FIG. 1, except for the provision of the high resistance layer 15. The high resistance layer 15 suppresses the current paths 12,13 of FIG. 1 that do not flow through the active region 6, and so improves the characteristics of the laser. The high resistance layer 15 may be, for example, a layer of silica (silicon dioxide).\nAlthough the high resistance layer 15 improves the characteristics of the laser, providing the high resistance layer 15 does make the process of fabricating the laser significantly more complicated. The high resistance layer must not be deposited over the contact layer 9, since this would increase the resistance of the desired current path 14 through the active region 6. Depositing the high resistance layer 15 therefore requires additional etching and masking steps—in general, the insulating layer 15 is deposited over the entire buffer layer 3 and the mesa, the portions of the insulating layer that are not over the mesa are masked, and the portion of the insulating layer that overlies the mesa is then removed via an etching step.\nIt is known to provide the laser structure with an additional layer that can be oxidised during the growth process to form the high resistance layer 15. However, if the additional layer extends through the mesa, it will affect the bandgap profile and refractive index profile of the laser, and this is undesirable. If the additional layer is to be deposited only outside the mesa, then additional masking and etching steps are required."}
{"text": "The present invention relates to safety devices, and more particularly relates to a vehicle safety device, such as for bicycles or the like, which is adapted to be detachably mounted adjacent the rear end thereof to provide a safety and warning signal for oncoming traffic.\nPresently, with the increased demand for and popularity of vehicles, such as bicycles, motorcycles, or the like, there has been a similar need for warning and/or signal devices for use therewith. Heretofor, various types of objects, such as reflectors, or flags, have been employed for this purpose. Such devices have generally included light colored reflectors, or flags, mounted by supporting wands or rods mounted adjacent the rear of the vehicle so as to not interfere with the rideroperator. For example, one such flag-type arrangement is disclosed in U.S. Pat. No. 3,947,059.\nHeretofor, such prior devices have not been completely satisfactory. Purely reflective type devices are generally helpful during night usage, but do not incorporate features which draw one's attention to same during daytime usage. Similarly, flag-type devices have the disadvantage, for example, that when becoming wet or moist, they have a tendency to stick, and hence, often do not provide the degree of movement required to attract the attention of oncoming traffic."}
{"text": "The present invention relates to a mechanism for adjusting tension of an inked ribbon of a printer.\nThe thermal printer which produces printed impressions by using an inked ribbon is well known.\nThere are two tension applying mechanisms, one of which is provided for applying a back tension to a feeding ribbon, and the other is provided for applying a winding-up tension. The value of the tension applied to the ribbon has influence on the quality of the printing.\nWhen the tension is too low, the ribbon wrinkles, causing printer failures in dots.\nIf the tension is too high, the ribbon slips and can not be fed.\nJapanese Patent Application Laid Open 7-89172 discloses a mechanism for controlling tension applied to an inked ribbon to a constant value by detecting the fluctuation of the tension of the inked ribbon which is caused by the change of diameter of the rolled ribbon during the printing operation.\nIn the conventional system, there must be provided a sensor for detecting the ribbon tension, and tension adjusting driving mechanisms in both of the ribbon feeding side and ribbon winding-up side. Consequently, the system becomes complicated in construction."}
{"text": "1. Field of Invention\nThe present invention relates to a plane bearing particularly used in an internal combustion engine.\n2. Description of Related Arts\nThe conventional plane bearing consists of a lining, a bonding layer formed on the lining, and an overlay. The lining consists of, for example, an aluminum bearing alloy. The bonding layer usually consists of Ni. The overlay consists of, for example, a Pb-based alloy which has good bearing property. The Pb-based overlay is described for example in DE 3000279AI.\nWhen the overlay wears out during sliding, the entire surface of the bonding layer is exposed and brought into contact with the shaft. Since the bearing property of the bonding layer are exceedingly inferior to those of the overlay and the lining, seizure due to adhesion with the shaft quickly occurs.\nThe present applicant has proposed in Japanese Unexamined Patent Publication No 62-283,216 that minute unevenesses are formed on the surface of the lining, and, further, the bonding layer and the overlay protrude into the concavities of lining. When the overlay of the proposed bearing wears out to a certain extent that top portions of the lining appear, the bonding layer and the overlay appear as islands surrounded by the top portions of the lining. Therefore seizure resistance is better than the conventional case where the overlay is worn out and only the bonding layer is exposed."}
{"text": "1. Field of the Invention\nThe present invention relates generally to hard-coated antiglare films, and polarizing plates and image displays including the same.\n2. Description of the Related Art\nWith technical improvements in recent years, liquid crystal displays (LCDs), plasma display panels (PDPs), electroluminescence displays (ELDs), etc. have been developed in addition to conventional cathode ray tube (CRT) displays as image displays and have been used in practical applications. Particularly, with technical innovations of LCDs with respect to, for example, wide viewing angles, high resolution, high response, and good color reproduction, applications of LCDs are spreading from laptop personal computers and monitors to television sets. In an LCD, generally, a liquid crystal panel is used in which polarizing plates are disposed on both sides of a liquid crystal cell, respectively. Generally, in order to prevent scratches on the polarizing plate, the surface of the liquid crystal panel is subjected to a hard coating treatment. For the hard coating treatment, a hard-coated film is often used. The hard-coated film is subjected to an antiglare treatment for preventing a decrease in contrast due to reflection of outside light such as fluorescent light and sunlight as well as reflection of an image at the surface of the liquid crystal panel. Particularly with an increase in size of the screens of image displays, the number of image displays equipped with hard-coated antiglare films is increasing.\nRecently, in order to improve image quality, the number of high definition image displays with small pixel sizes is increasing. In such high definition image displays, when a conventional hard-coated antiglare film is disposed, variations in brightness in pixels are emphasized further more to cause a visible failure (a failure due to glare) to result in considerably deteriorated image quality. Conventionally, in order to produce an antiglare laminate that allows high definition to be provided, a method has been employed in which the haze value of the antiglare layer is increased to eliminate glare. However, this method has a problem in that the contrast is decreased considerably because light is scattered intensively at the panel surface. Furthermore, in the case where the surface unevenness is increased to improve antiglare properties, there is a problem in so-called white blur in an oblique direction, that is, white blur is observed due to excessively intensive scattering of reflected light when the panel surface is viewed from the oblique direction. In the antiglare treatment, a method of producing an uneven shape of the film surface by adding, for example, inorganic or organic particles is used. Generally, although an improvement in antiglare properties and an improvement in contrast or white blur are considered to have a contradictory relationship, various proposals have been made to obtain both these properties. For example, it is being studied that aggregates with three-dimensional structures formed of the aforementioned particles are allowed to be present in an antiglare layer (see, for instance, JP 2005-316413 A). This, however, may cause scatter at the aggregates or may cause a fine pattern to appear in a hard-coated film. Furthermore, although a method that is effective in improving a part of properties has been proposed (see, for instance, JP 2003-4903 A and JP 4001320 B), a method that is effective in solving all the three problems described above has not been found."}
{"text": "The delivery of carefully-planned doses of radiation may be used to treat various medical conditions. For example, radiation treatments are used, often in conjunction with other treatments, in the treatment and control of certain cancers. While it can be beneficial to deliver appropriate amounts of radiation to certain structures or tissues, in general, radiation can harm living tissue. It is desirable to target radiation on a target volume containing the structures or tissues to be irradiated while minimizing the dose of radiation delivered to surrounding tissues. Intensity modulated radiation therapy (IMRT) is one method that has been used to deliver radiation to target volumes in living subjects.\nIMRT typically involves delivering shaped radiation beams from a few different directions. The radiation beams are typically delivered in sequence. The radiation beams each contribute to the desired dose in the target volume.\nA typical radiation delivery apparatus has a source of radiation, such as a linear accelerator, and a rotatable gantry. The gantry can be rotated to cause a radiation beam to be incident on a subject from various different angles. The shape of the incident radiation beam can be modified by a multi-leaf collimator (MLC). A MLC has a number of leaves which are mostly opaque to radiation. The MLC leaves define an aperture through which radiation can propagate. The positions of the leaves can be adjusted to change the shape of the aperture and to thereby shape the radiation beam that propagates through the MLC. The MLC may also be rotatable to different angles.\nObjectives associated with radiation treatment for a subject typically specify a three-dimensional distribution of radiation dose that it is desired to deliver to a target region within the subject. The desired dose distribution typically specifies dose values for voxels located within the target. Ideally, no radiation would be delivered to tissues outside of the target region. In practice, however, objectives associated with radiation treatment may involve specifying a maximum acceptable dose that may be delivered to tissues outside of the target.\nTreatment planning involves identifying an optimal (or at least acceptable) set of parameters for delivering radiation to a particular treatment volume. Treatment planning is not a trivial problem. The problem that treatment planning seeks to solve involves a wide range of variables including:                the three-dimensional configuration of the treatment volume;        the desired dose distribution within the treatment volume;        the locations and radiation tolerance of tissues surrounding the treatment volume; and        constraints imposed by the design of the radiation delivery apparatus.        \nThe possible solutions also involve a large number of variables including:                the number of beam directions to use;        the direction of each beam;        the shape of each beam; and        the amount of radiation delivered in each beam.        \nVarious conventional methods of treatment planning are described in:                S. V. Spirou and C.-S. Chui. A gradient inverse planning algorithm with dose-volume constraints, Med. Phys. 25, 321-333 (1998);        Q. Wu and R. Mohand. Algorithm and functionality of an intensity modulated radiotherapy optimization system, Med. Phys. 27, 701-711 (2000);        S. V. Spirou and C.-S. Chui. Generation of arbitrary intensity profiles by dynamic jaws or multileaf collimators, Med. Phys. 21, 1031-1041 (1994);        P. Xia and L. J. Verhey. Multileaf collimator leaf sequencing algorithm for intensity modulated beams with multiple static segments, Med. Phys. 25, 1424-1434 (1998); and        K. Otto and B. G. Clark. Enhancement of IMRT delivery through MLC rotation,” Phys. Med. Biol. 47, 3997-4017 (2002).        \nAcquiring sophisticated modern radiation treatment apparatus, such as a linear accelerator, can involve significant capital cost. Therefore it is desirable to make efficient use of such apparatus. All other factors being equal, a radiation treatment plan that permits a desired distribution of radiation dose to be delivered in a shorter time is preferable to a radiation treatment plan that requires a longer time to deliver. A treatment plan that can be delivered in a shorter time permits more efficient use of the radiation treatment apparatus. A shorter treatment plan also reduces the risk that a subject will move during delivery of the radiation in a manner that may significantly impact the accuracy of the delivered dose.\nDespite the advances that have been made in the field of radiation therapy, there remains a need for radiation treatment methods and apparatus and radiation treatment planning methods and apparatus that provide improved control over the delivery of radiation, especially to complicated target volumes. There also remains a need for such methods and apparatus that can deliver desired dose distributions relatively quickly."}
{"text": "The invention generally relates to using cameras in connection with a marine seismic survey.\nSeismic exploration involves surveying subterranean geological formations for hydrocarbon deposits. A survey typically involves deploying seismic source(s) and seismic sensors at predetermined locations. The sources generate seismic waves, which propagate into the geological formations creating pressure changes and vibrations along their way. Changes in elastic properties of the geological formation scatter the seismic waves, changing their direction of propagation and other properties. Part of the energy emitted by the sources reaches the seismic sensors. Some seismic sensors are sensitive to pressure changes (hydrophones), others to particle motion (e.g., geophones), and industrial surveys may deploy only one type of sensors or both. In response to the detected seismic events, the sensors generate electrical signals to produce seismic data. Analysis of the seismic data can then indicate the presence or absence of probable locations of hydrocarbon deposits.\nSome surveys are known as “marine” surveys because they are conducted in marine environments. However, “marine” surveys may be conducted not only in saltwater environments, but also in fresh and brackish waters. In one type of marine survey, called a “towed-array” survey, an array of seismic sensor-containing streamers and sources is towed behind a survey vessel."}
{"text": "The predominate approach today to introduce factory automated technology into manufacturing is to selectively apply automation and to create islands of automation. The phrase \"islands of automation\" has been used to describe the transition from conventional or mechanical manufacturing to the automated factory. Interestingly, some appear to use the phrase as though it were a worthy end object. On the contrary, the creation of such islands can be a major impediment to achieving an integrated factory.\nManufacturing examples of islands of automation often include numerically controlled machine tools; robots for assembly, inspection, painting, and welding; lasers for cutting, welding and finishing; sensors for test and inspection; automated storage/retrieval systems (AS/RS) for storing work-in-process, tooling and supplies; smart carts, monorails, and conveyors for moving material from work station to work station; automated assembly equipment and flexible machining systems. Such islands are often purchased one at a time and justified economically by cost reductions. An example of an AS/RS system is disclosed in the U.S. Pat. No. 4,328,422 to Loomer. A different type of AS/RS system and control system therefor is disclosed in the U.S. Pat. No. 4,232,370 to Tapley.\nTo integrate the islands of automation it is necessary to link several machines together as a unit. For example, a machine center with robots for parts loading and unloading can best be tied to visual inspection systems for quality control Computer numerical control machine tools can all be controlled by a computer that also schedules, dispatches, and collects data. Selecting which islands to link can be most efficiently pursued on the basis of cost, quality and cycle time benefits.\nIn some cases the islands of automation will be very small (e.g. an individual machine or work station). In other cases the islands might be department-sized. The U.S. Pat. No. 4,611,749 to Kawano discloses the use of robots to transfer parts between such islands which are relatively close to each other.\nFrom a systems viewpoint, islands of automation are not necessarily bad, so long as they are considered to be interim objectives in a phased implementation of an automated system. However, to obtain an integrated factory system, the islands of automation must be tied together or synchronized. Systems synchronization frequently occurs by way of a material-handling system; it physically builds bridges that join together the islands of automation. Early examples of such islands of automation linked together by a material-handling system are disclosed in the U.S. Pat. Nos. 4,369,563 to Williamson and 3,854,889 to Lemelson.\nThe '563 patent discloses a system including machine tools which perform machining operations on workpieces loaded on pallets. The pallets are delivered to the machine tools from a storage rack by transporters. The workpieces are manually loaded onto the pallets.\nThe '889 patent discloses a system including work-holding carriers which are selectively controlled in their movement to permit work to be transferred to selected machine tools while bypassing other machine tools.\nAutomated material handling has been called the backbone of the automated factory. Other than the computer itself, this function is considered by many automation specialists as the most important element in the entire scenario of automated manufacturing. It is the common link that binds together machines, workcells, and departments into a cohesive whole in the transformation of materials and components into finished products. For example, the U.S. Pat. No. 4,332,012 to Sekine et al discloses a control system for assembly lines for the manufacture of different models of automotive vehicles. Temporary storage is provided between assembly steps by a storage section.\nTo date, the major application for industrial robots has been material handling. Included here are such tasks as machine loading and unloading; palletizing/depalletizing; stacking/unstacking; and general transfer of parts and materials--for example, between machines or between machines and conveyors. An example of one such application is disclosed in the U.S. Pat. No. 4,519,761 to Kenmochi. The '761 Patent discloses a combined molding and assembling apparatus wherein a pallet is conveyed by a conveyor. Resin components are carried by the pallet for use in the molding and assembling operation.\nRobots are often an essential ingredient in the implementation of Flexible Manufacturing Systems (FMS) and the automated factory. Early examples of the use of robots for assembling small parts is disclosed in the U.S. Pat. Nos. 4,163,183 and 4,275,986 wherein robots are utilized to assemble parts from pallets onto a centrally located worktable.\nThe automated factory may include a variety of material transportation devices, ranging from driver-operated forklifts to sophisticated, computer-operated, real-time reporting with car-on-track track systems and color graphics tracking. These material transport systems serve to integrate workcells into FMS installations and to tie such installations and other workcells together for total factory material transport control.\nWith all of their versatility, robots suffer from a limitation imposed by the relatively small size of their work envelope, requiring that part work fixtures and work-in-process be brought to the robot for processing. Complete integration of the robot into the flexible manufacturing system requires that many parts and subassembIies be presented to the robot on an automated transport and interface system. For example, installation of an assembly robot without an automated transport system will result in an inefficient island of automation needing large stores of work-in-process inventory for support, which are necessary to compensate for the inefficiencies of manual and fork truck delivery.\nA recent example of the use of robots in a manufacturing assembly line is disclosed in the U.S. Pat. No. 4,611,380 to Abe et al. The '380 Patent also discloses the use of a bar code to identify the components to be assembled to a base component to control the assembly operations.\nThe U.S. Pat. No. 4,616,411 to Suzuki et al discloses a fastening apparatus including a bolt receiving and supply device for use in the automated assembly of a door to a vehicle.\nThe handling, orienting and feeding of parts as they arrive from vendors are formidable jobs which must be done prior to robotic assembly since, in general, all such parts require reorienting for the assembly robot. The U.S. Pat. No. 4,527,326, to Kohno et al., for example, discloses a vibratory bowl which feeds parts to an assembly robot. A vision system enables the robot to properly pick up the parts from the bowl.\nPart feeding is a technology that generally has lagged behind the advanced automation systems it supports. However, in general, part feeding curtails flexibility, increases costs, increases floor space requirements and lengthens concept-to-delivery time. For maximum flexibility, a minimum amount of tooling should be considered. On the other hand, additional tooling can be used effectively to \"buy time\" by assisting the robot. Typically, dedicated hardware--bowl feeders, magazines, pallets--is required to feed parts to the robot. Unlike the robot, dedicated hardware is not easily reusable and therefore is less economical for medium-volume applications.\nThe U.S. Pat. No. 4,383,359 discloses a part feeding and assembly system, including multiple stage vibration and magazine feeders. A robot is utilized to change the position of the fed parts for assembly on a chassis supported on a line conveyor. The robot operates in combination with a vision system to reorient the parts.\nNeither flexible nor sophisticated, part feeding equipment is usually constructed by highly skilled artisians working with welding torch and hammer in small specialized shops. The most common and most inexpensive feeding method--vibratory bowl feeding--provides the builder with a versatile base easily modified to handle many different parts which are not delicate and which are substantially identical. Delicate parts or parts that tangle, such as motors, are better fed by magazines or trays for exact orientation.\nAlso, not all parts, for example, can be bowl fed. For most parts, the overriding concern is geometry and, in particular, symmetry. If a part is either symmetric or grossly asymmetric, then vibratory bowl feeding will be easier and more efficient.\nRobots may load and unload workpieces, assemble them on the transport, inspect them in place or simply identify them. The kind of activity at the robot or machine and material transport system interface dictates the transport system design requirements. One of the design variables relating to the interface includes accuracy and repeatability of load positioning (in three planes). Also, care in orienting the workpiece when it is initially loaded onto the transport carrier will save time when the work is presented to the robot or the tool for processing. Proper orientation of the part permits automatic devices to find the part quickly without \"looking\" for it and wasting time each time it appears at the workstation.\nFixtures may be capable of holding different workpieces, reducing the investment required in tooling when processing more than one product or product style on the same system.\nThe transport system must be capable of working within the space limitations imposed by building and machinery configurations, yet must be capable of continuous operation with the loads applied by a combination of workpiece weight, fixture weight, and additional forces imposed by other equipment used in the process.\nThe system must also have the ability to provide queuing of parts at the workstation so that a continuous flow of work is maintained through the process. Automatic queuing of transport carriers should provide gentle accumulation without part or carrier damage.\nThe primary impediment to robotic assembly is economic justification. When the cost of robotic assembly is compared against traditional manual methods or high volume dedicated machinery, robots oftentimes lose out. On one side of the spectrum are the high-volume, high-speed applications where hard automation is used. It's difficult for robots to compete in that environment. On the other side are the low-volume, high variety products that are assembled manually. Robots may lack the dexterity to perform these jobs, and they may cost more than relatively low-paid manual assemblers. There is a middle ground between these two extremes for flexible assembly. Many believe that the best approach is a combination of robots, dedicated equipment and manual assembly.\nThere are other barriers to the use of robots in mechanical assembly. They include the following: (1) the high cost of engineering a new system, which may run three to five times the cost of the robot itself; (2) the amount of time it takes to engineer the system; (3) the difficulty of coordinating multiple arms; (4) the difficulty of integrating an assembly system; (5) the high cost of tooling, software sensors, part presentation equipment, and other peripherals; (6) the difficulty of finding knowledgeable personnel; (7) insufficient speed, lift capacity, and positioning accuracy and repeatability on the part of the robots; and (8) a lack of supporting technology in such areas as high-level programming languages, end-of-arm tooling, and sensors.\nAdditional impediments to the successful implementation of robots in assembly are insufficient communication among departments, a general slowdown in capital equipment acquisition, a disinclination to plan ahead, fear of change, and the infamous NIH (Not Invented Here) Syndrome that keeps companies from accepting ideas originating outside their walls.\nStill, while assembly is probably the most difficult area of robotic application, many say it also holds the most promise. Assembly robots offer an array of benefits that cannot be ignored. They can produce products of high and consistent quality, in part because they demand top-quality components. Their reprogrammability allows them to adapt easily to design changes and to different product styles. Work-in-process inventories and scrap can be reduced. Therefore, it is important that the materials transport system serving the robots be capable of quickly moving into position with parts, then quickly moving out of the workstation and on to downstream stations. Prompt transporter movements between stations allow work-in-process inventory to be minimized. Batch sizes are smaller and work faster with only a minimum of queuing at each workstation.\nThe U.S. Pat. No. 4,594,764 to Yamamoto discloses an automatic apparatus and method for assembling parts in a structure member such as an instrument panel of an automobile. A conveyor conveys a jig which supports the panel to and from assembly stations. Robots mount the parts on the instrument panel at the assembly stations. Robots are provided with arm-mounted, nut-driving mechanisms supplied from vibratory parts bowls.\nA link for tying together some of the independently automated manufacturing operations is the automatic guided vehicle system (AGVS). The AGVS is a relatively fast and reliable method for transporting materials, parts or equipment, especially when material must be moved from the same point of origin to other common points of destination. Guide path flexibility and independent, distributed control make an AGVS an efficient means of horizontal transportation. As long as there is idle space and a relatively smooth floor to stick guide wires or transmitters into, the AGVS can be made to go there.\nAs an alternative to traditional conveying methods, the AGVS provides manufacturing management with a centralized control capability over material movement. Also, the AGVS occupies little space compared with a conveyor line. Information available from the AGVS also provides management with a production monitoring data base. The U.S. Pat. No. 4,530,056 to Mackinnon et al. discloses an AGVS system including a control system for controlling the individual vehicles.\nA relatively new type of link for tying together independently automated manufacturing operations is the mechanical guided vehicle (MGV) commercially available from the Roberts Corporation of the Cross and Trecker Corporation of Bloomfield Hills, Mich. The MGV is a self-contained, self-propelled, battery-operated vehicle which travels on a track system. The vehicle is utilized to carry a load much like an automated guided vehicle.\nRobot installations for transporter interface can be grouped into three principal categories: (1) stationary robots, (2) moving (i.e. mobile) robots (on the floor or overhead), and (3) robots integral with a machine. The moving robots subdivide into two types. First are stationary robots, mounted on a transporter to move between work positions to perform welding, inspection, and other tasks. The second type of moving robot is the gantry unit that can position workpieces weighing more than one ton above the workcells and transport system. The system only has to deliver and pick up somewhere under the span of gantry movement.\nEnd effectors used in material handling include all of the conventional styles--standard grippers, vacuum cups, electromagnets--and many special designs to accommodate unusual application requirements. Dual-purpose tooling is often used to pick up separators or trays, as well as the parts being moved through the system.\nVacuum-type grippers and electromagnetic grippers are advantageous because they permit part acquisition from above rather than from the side. This avoids the clearance and spacing considerations that are often involved when using mechanical grippers.\nHowever, the use of vacuum and electromagnetic grippers is not without its problems since cycle time is not just a function of robot speed and its accelerating/decelerating characteristics. Cycle time is dependent on how fast the robot can move without losing control of the load. Horizontal shear forces must be considered in the application of these grippers. This often means that the robot is run at something less than its top speed."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a semiconductor device.\n2. Description of the Related Art\nIn recent years, efforts have been made to develop a ferroelectric memory that uses a ferroelectric film as a dielectric film for a capacitor. In a ferroelectric memory, noble metal such as platinum is used for capacitor electrodes in order to suppress reaction of the ferroelectric film during crystallization annealing (see, for example, Jpn. Pat. Appln. KOKAI Publication No. 9-162311).\nHowever, a product (noble metal compound) obtained by etching noble metal such as platinum exerts a low saturated vapor pressure. This makes it difficult to appropriately achieve anisotropic etching. Accordingly, appropriate anisotropic etching requires substrate temperature to be increased to about 300 to 400° C. It is thus difficult to use a photoresist as an etching mask. Consequently, anisotropic etching is generally carried out using a hard mask such as a silicon oxide film.\nA hard mask is formed by a film forming method such as CVD or sputtering. A film formed by such a film forming method generally adheres well. Consequently, the hard mask cannot be easily removed by the above method.\nOn the other hand, the reduced size of semiconductor devices has made it difficult to precisely process gate electrodes of MIS transistors by anisotropic etching. To form gate electrodes by anisotropic etching, it is necessary to maximize the etching selective ratio of gate electrode films to gate insulating films.\nHowever, a higher etching selective ratio enlarges the lower part of the gate electrode, which thus has an inappropriate shape. An attempt to reduce the etching selective ratio to improve the gate electrode shape may disadvantageously result in etching the gate insulating film, exposing a semiconductor substrate.\nAs described above, if an adhesive film such as a hard mask is used as an etching mask to etch an etching target film, removing the etching mask is disadvantageously difficult.\nFurther, if the etching target film such as the gate electrode film is etched, it is disadvantageously difficult to obtain an appropriate etching processing shape without affecting an underlying film such as the gate insulating film."}
{"text": "1. Field of the Invention\nThis invention relates to textiles containing polyester fibers and particularly relates to compositions which can impart soil-release properties to polyester fibers after application to the textile and curing. It specifically relates to compositions which are partially esterified interpolymers of maleic anhydride for use as soil release finishes.\n2. Description of the Prior Art\nSynthetic fiber-containing clothing has brought great changes in the textile industry, partly because these clothes can retain an initial crease and partly because they can be washed with water. Additional crease-retention is generally imparted by treating the textile, before manufacture of the clothing, with resins and curing the resin on the textile after creasing and pressing as desired, thereby providing the so-called \"wash-and-wear\" and \"durable press\" clothing so plentiful in the clothing industry. However, clothing articles which contain polyester fibers also have the disadvantage of more readily picking up soil, and especially oily soil, than does natural-textile clothing. Moreover, when washed in a home washing machine, the polyester-containing clothing may become progressively more heavily soil-laden with successive washings because soil which is dispersed in the wash water redeposits onto the washed clothing before the end of the wash cycle.\nThe oleophilic and hydrophobic properties of linear polyester fibers, which are normally blended with cellulosics such as cotton and rayon, are responsible for this redeposition behavior. These polyester fibers actually seem to scavenge oily soil from the wash water. Aminoplast resins used in durable press finishing also have a tendency to retain oil-and water-borne stains. The combination of fairly high resin loads and high polyester content creates a serious soil retention problem for durable-press garments and other textile articles.\nThis behavior of polyester-containing textiles and several attempts to solve the problem are discussed in considerable detail in U.S. Pat. No. 3,377,249. In general, these prior attempts to improve the soil-resisting properties of polyester-containing cellulosic fabrics have used nitrogen-containing resins to impart hydrophilic qualities to the polyester fibers. The above-mentioned patent also discloses the application of aminoplast resins, in combination with a catalyst and a synthetic acid-emulsion polymer based on arcylic acid, to a textile which contains linear polyester fibers, whereby soil-release and durable-press characteristics are imparted thereto. However, these new developments have not completely solved the problem.\nSeveral textile finishing compositions which counteract the soil-retention tendency of polyester-containing fabrics are now on the market. The two leading commercial types of soil-release textile finishes are the acrylic-emulsion type, such as Rhoplex SR-488, and the fluorocarbon-based type, such as 3M Scotchgard. The composition of the instant invention forms a third and advantageous type."}
{"text": "1. Field of the Invention\nThe present invention relates, generally, to the field of investment management and, more particularly, to a system, method, and computer program product for allocating assets among a plurality of investments to guarantee a predetermined value at the end of a predetermined time period.\n2. Discussion of the Background\nPresently, there is a wide array of investments that investors may use to achieve their investment goals. As is well known in the investment industry, each investment has its own advantages and disadvantages, including an associated anticipated rate of return and risk. Typically, an investment with a high anticipated rate of return also has a high risk of depreciating. For example, while some stocks may provide very high rates of return, they also involve substantial risk that the investor could lose part or all of the initial investment. Well-known examples of high risk investments include stocks (especially technology stocks) and mutual funds (especially growth-oriented funds). Conversely, an investment that is low risk, typically has an anticipated rate of return that is relatively low. Examples of lower risk investments include insured savings accounts, Certificate of Deposits (CDs), government backed bonds, and money market accounts.\nTraditionally, investors have had to choose between a higher rate of return (selecting a high yield, high risk investment) and safety of principal (selecting a low risk, low yield investment). In seeking a moderately high rate of return while attempting to protect against significant downturns in the market, investors often diversify—allocating their assets among higher risk, higher yield investments (hereinafter referred to as “high yield investments”) and lower risk, lower yield investments (hereinafter referred to as “low risk investments”). The investor's assets are allocated between the high yield and low risk investments according to the investor's risk tolerance and investment goals.\nWhile diversifying in this fashion may lower the risk of significant downturns in the value of the investment portfolio, it also tends to reduce the potential for appreciation. Diversification tends to lower the downside risk by reducing the depreciation that may result in the total value of the investor's assets if the high yield investment yields a negative return (i.e., experiences a reduction in value). However, diversifying can also result in limited potential for appreciation. Specifically, if the high yield investment produces a high rate of return on the assets allocated therein—as is anticipated and intended—the increase in the overall value of the investor's assets will be diminished by the relatively low increase in value of the low risk investment. For example, if half of an investor's assets are allocated to a high yield investment that yields a fifty percent rate of return, and the other half of the investor's assets are allocated to a low risk investment that yields a zero percent return, the overall rate of return on the investor's assets will be twenty-five percent. Thus, while diversifying may increase the overall security of the investor's investment portfolio, diversifying may also reduce the overall rate of return.\nFinancial institutions have also created various investment vehicles designed to reduce or eliminate risk, while minimizing the performance dampening effect of the diversification. However, these vehicles typically place restrictions on the investment and/or investor and also include other undesirable features.\nZero coupon instruments, certificates of deposit, guaranteed investment contracts and variations on these types of instruments comprise one type of such an investment vehicle. With this type of instrument, a financial institution guarantees that if the principal contribution is maintained without withdrawals or transfers until the specified date, the investment will be worth a predetermined value at the end of the period. The financial institution invests a pool of investments standing behind the obligations to every investor, be it individuals or institutions, as it sees fit, subject to applicable regulatory standards, to support its guarantee. Typically, these types of investment do not yield greater returns than the predetermined value at the end of the predetermined time period and are often a fixed rate investment. In addition, the investor usually has no say in the selection of investments and, once the assets are allocated, the investor has no control to reallocate them.\nIndividualized guaranteed portfolios are another type of such program. With this type of instrument, a financial institution negotiates with an investor—generally an institutional investor—over the range of investments that may be used in a portfolio underlying a guarantee of some predetermined value at the end of the predetermined period. In such a negotiation, the predetermined value and the length of the period, as well as the ratings of the issuers of the underlying securities, may all be part of the negotiation. If the predetermined value is not accumulated at the end of the predetermined time period, the financial institution supplies the shortfall. However, the financial institution taking on the guarantee then maintains control over the investment of the portfolio once the period begins, within the parameters set out in the negotiation. Thus, the investor has limited selection of investments and, once the assets are allocated, the investor has no control to reallocate them.\nA guaranteed return pooled investment is still another type of such program. These instruments generally are offered by investment companies, and typically comprise a particular investment such as a mutual fund with a limited “offering period” (the period within which an investment must be made) and a set calendar date, which is the date the pre-determined value is guaranteed for each shareholder. The investor has no control over the selection of investments and once the assets are allocated, the investor has no control to reallocate them. In summary, the prior art products impose restrictions on the investments available for initial allocations, restrict or prohibit reallocations, and often have a relatively low anticipated rate of return.\nTherefore, notwithstanding the available investment diversification products, there is a need for a system, method, and computer program product for allocating assets among a plurality of investments that, unlike such conventional diversification products: (1) guarantees a predetermined value or rate of return on an investment at the end of a predetermined time period; (2) has a guaranteed rate of return over a predetermined time period, without the reduction in the appreciation normally associated with diversification; (3) allows for diversification between a Secure Account and a Variable Account to guarantee a predetermined value at the end of a predetermined time period and to increase the investor's overall security of the portfolio; (4) provides significant investment control to the investor with flexibility in the type of investments available to the investor throughout the time period of the investment; and (5) does not have a fixed asset allocation requirement."}
{"text": "Types of non-volatile memory such as NAND memory typically suffer from a limited number of program/erase cycle counts before these types of memory perform unacceptably and/or become unreliable. For example, some types of NAND memory store charge on a floating gate that is surrounded by an insulating oxide layer to keep the charges trapped. NAND memory operations depend on the quality of the oxide layer between a memory cell channel and the floating gate. As program/erase cycle counts increase the oxide layer progressively wears down or degrades due to charges (electrons/holes) flowing through the oxide layer at each program/erase cycle. This oxide layer degradation may result in slower erase operations, higher bit-error rates (BER) during programming and charge leakage from the floating gate. All three of these symptoms of oxide layer degradation may limit NAND memory endurance.\nNAND memory endurance may also be limited by charge trapped in an oxide above a lightly doped drain (LDD) region of a NAND memory cell. Trapped charges may cause an apparent increase in resistance of a memory string (bitline). In a worst case, no current passes through the memory string and the memory string may appear as an open circuit. Appearing as an open circuit may lead to all bits on that memory string to read as programmed “0” even after erase operations. As a result of always reading as programmed, these memory cells will no longer be used."}
{"text": "Apparatus and methods are known for accessing a patient's vasculature percutaneously to perform a procedure within the vasculature, and for sealing the puncture that results after completing the procedure. For example, a hollow needle may be inserted through a patient's skin and overlying tissue into a blood vessel. A guide wire may be passed through the needle lumen into the blood vessel, whereupon the needle may be removed. An introducer sheath may then be advanced over the guide wire into the vessel, e.g., in conjunction with or subsequent to one or more dilators. A catheter or other device may be advanced through the introducer sheath and over the guide wire into a position for performing a medical procedure. In this manner, the introducer sheath may facilitate introducing various instruments into the vessel, while minimizing trauma to the vessel wall and/or minimizing blood loss. Upon completing the procedure, the instrument(s) and introducer sheath may be removed, leaving a puncture extending between the skin and the vessel.\nTo seal the puncture, external pressure may be applied to the overlying tissue, e.g., manually and/or using sandbags, until hemostasis occurs. This procedure, however, may be time consuming and expensive, requiring as much as an hour of a medical professional's time. It is also uncomfortable for the patient, and may require the patient to remain immobilized in an operating room, catheter lab, or holding area. In addition, a risk of hematoma exists from bleeding before hemostasis occurs.\nVarious apparatus and methods have been suggested for sealing a percutaneous puncture instead of using external pressure. For example, U.S. Pat. No. 5,108,421 to Fowler discloses a collagen plug that may be delivered into a puncture through tissue. In one embodiment, a catheter is inserted through the puncture into the blood vessel. A balloon on the catheter is expanded and retracted until the balloon is disposed adjacent the puncture at the wall of the vessel. The plug may be advanced into the puncture until the plug contacts the balloon, thereby preventing the plug from entering the vessel. Once the plug is positioned within the puncture, the balloon may be deflated and withdrawn, leaving the plug therein to expand and seal the puncture and/or to promote hemostasis.\nBy way of another example, U.S. Pat. Nos. 5,192,302 and 5,222,974 issued to Kensey et al. describe a biodegradable collagen plug that may be delivered through an introducer sheath into a puncture site. The disclosed plug, however, may be difficult to position properly with respect to the vessel, which may be significant since it is generally undesirable to expose the collagen material within the bloodstream where it may float downstream and cause an embolism."}
{"text": "High-speed digital data links may suffer from inter-symbol interference, especially in situations in which loss, reflections or other imperfections exist in the transmission channel. Inter-symbol interference may have the effect that the signal received during a given clock cycle is a linear combination of the bit transmitted during the corresponding clock cycle at the transmitter, and of the bits transmitted during a number of preceding clock cycles. The effects of inter-symbol interference may be mitigated using a technique referred to as decision feedback equalization (DFE) which involves correcting the received signal at the sampling point, during each clock cycle, with a linear combination of the bits received during a number of preceding clock cycles.\nThe contribution from the immediately preceding received bit, which is referred to as the first tap, may be generated using a technique referred to as predictive decision feedback equalization (predictive DFE, which may also be referred to as speculative DFE or loop-unrolled DFE), in which two correction terms are calculated, one corresponding to a received 1, and one corresponding to a received 0 in the immediately preceding received bit; the appropriate one of these two correction terms is then selected using a multiplexer (MUX) once a binary value for the bit received on the immediately preceding clock cycle is available. The processing of the correction from the last received bit may impose relatively stringent requirements on the operating speed of the circuit.\nThus, there is a need for a system for predictive DFE with relaxed timing requirements for circuits in the DFE loop."}
{"text": "1. Field of the Invention\nThe invention relates to a method and an apparatus for matching the flow capacity of a radial turbine of a turbocharger to a flow capacity of an internal combustion engine.\n2. Discussion of Background\nRadial turbines for turbochargers typically include at the gas inlet end, of a spiral-shaped gas inlet casing and a nozzle ring spatially terminating this casing in the direction toward the gas outlet end. The rotor, which is guided on the compressor drive shaft and is provided with vanes, is arranged at the gas outlet end. The gas flow duct extends between the rotor hub and the opposite region of the gas outlet casing. Instead of a gas outlet casing, it is also known to use a gas outlet flange which is releasably connected to the gas inlet casing by a plurality of bolts uniformly distributed on its periphery.\nIf a turbocharger is to be associated with a certain internal combustion engine, the swallowing capacity of the radial turbine, i.e. the size of the delivery flow processed by it for a given machine diameter, has to be correspondingly matched.\nFor this purpose, a method is known in which the rotor vanes and the nozzle ring vanes are first shortened by removal of material from the vane height to correspond with the parameters of the internal combustion engine. As an alternative to this and--if required--additionally also, the angle of incidence of the nozzle ring vanes can or must be modified, i.e. different nozzle rings with the required vane angle of incidence are employed. Because the gas inlet casing and the nozzle ring are accurately matched to one another, the gas inlet casing must subsequently be correspondingly matched, i.e. it must also be exchanged.\nIn this method, the matching of the radial turbine to correspond with the modified rotor takes place at the gas inlet end, i.e. both the nozzle ring and the spiral-shaped gas inlet casing must be exchanged. For this reason, the method is very expensive in terms of labor and is therefore costly. Because of the many combination possibilities, this requires quite a high stores holding of the corresponding exchange parts and very extensive transport requirements for the matching of the turbocharger to the corresponding internal combustion engine (which takes place on site). This further increases the costs.\nIn special cases where a turbocharger is matched to an internal combustion engine, such as in emergency electricity generating units, where a rapidly responding turbocharger is required, the mass moment of inertia of the rotor is reduced by employing an appropriate rotor in a lower mass material, for example, ceramic. These known rotors, however, are very expensive and, furthermore, do not have sufficient functional reliability."}
{"text": "1. Field of the Invention\nThe present invention is directed to a flow transducer muff for positioning a substantially right-angled parallelepiped-shaped flow transducer in the expiration channel of a ventilator, of the type having first and second muff halves which are connectable on opposite sides of the flow transducer.\n2. Description of the Prior Art\nExpired gas flow is measured in a flow transducer at the expiration side of a ventilator. The flow transducer is maintained in position in the expiration channel by means of a flow transducer muff, which forms a gas-tight connection to the flow transducer. Such a flow transducer muff is described, for example, in the Operating Manual (page 10: 2-3, August 1988) for the Servo Ventilator 900C, sold by Siemens-Elema AB. In general, known flow transducer muffs consist of two halves which are connected on opposite sides of the flow transducer.\nThe two muff halves in known components are identical in shape, and there is thus a risk that the flow transducer could be installed, after cleaning, in a reversed position in the ventilator."}
{"text": "From DE 198 21 164 A1 has become known a typical six-gear dual-clutch transmission in which two fixed wheels are situated upon a first input shaft designed as hollow shaft. While the first fixed wheel is in tooth contact with the idler wheel for the reverse gear and a second idler wheel for the second gear, the second fixed wheel meshes with one idler wheel for the sixth gear and with one idler wheel for the fourth gear. In addition, the idler wheels for the reverse gear and the sixth gear are supported in the transmission upon a first countershaft while the idler wheels for the second gear and the fourth gear are placed upon a second countershaft. This dual-clutch transmission has a comparatively short construction so that it can be used for crosswise installation in a motor vehicle.\nFrom DE 199 23 185 A1 has further become known a dual-clutch transmission in which both, upon a hollow first input shaft and upon the second input shaft, are supported therein and designed as solid shaft, a respective fixed wheel is situated which drives each of two idler wheels of different transmission gears. The first fixed wheel situated upon the hollow shaft thus meshes with the idler wheel for the third gear and the idler wheel for the fifth gear. The second fixed wheel situated upon the solid shaft is, on the other hand, in tooth contact with the idler wheels for the second and for the fourth gears.\nWith this background, the problem on which the invention is based is to introduce a structure for dual-clutch transmissions which can be used both in a six-gear transmission and in a seven-gear transmission. Besides, a six-gear transmission thus built must be shorter than the known generic transmission drafts and the seven-gear dual-clutch transmission must fit in a similar vehicle installation space which hitherto had been adequate only for dual-clutch transmission having less gears. Finally, by reducing the hitherto needed transmission parts, the production cost is lowered.\nThe solution of this problem results from the features of the main claims while advantageous developments and improvements of the invention can be understood form the sub-claims."}
{"text": "In complex communication systems, such as traditional telephone, voice over Internet protocol (VoIP), session initiation protocol (SIP), and/or the like, messages typically must pass through a variable number of hops, in order to reach their final destination. For example, in a SIP environment, the total time (end-to-end time) that it takes a message to be processed through the plurality of hops often determines as to whether or not the message is ultimately successful in reaching its final destination through the environment. In order to attempt to manage the volume of messages through their environment so they are successfully transmitted, communication system carriers typically allocate a fixed time for each hop of the network (or, system). In this manner, the overall time for the furthest (or, maximum) path (i.e., largest quantity of hops) does not exceed the allotted time, and the overall time equals merely the quantity of hops multiplied against this “per hop” fixed allocated time. This rigid and crude management of messages through the system does not account for message-specific items and is both inefficient and inflexible. For example, there is no prioritization of messages based on any other factors beyond merely allocating the fixed per hop time.\nIn view of the foregoing, a need exists to overcome one or more of the deficiencies in the related art."}
{"text": "This invention relates generally to lifting devices and, more particularly, to a lifting device that attaches to the front or rear of a vehicle for lifting the front or rear wheels off the ground.\nWhen needing to change a tire, install tire chains, or needing to work on the underside of the vehicle, it is sometimes difficult to find a good place along the side of the vehicle to position a jack. Another problem with vehicle jacks is that they are sometimes insecure due to uneven ground, lack of a stable position beneath the vehicle frame, or jack product failure.\nVarious devices have been proposed for a vehicle jack used in cooperation with a vehicle receiver hitch. While the existing devices and proposals are assumably effective for their intended uses, the existing devices do not enable a user to effectively and safely raise either end of a vehicle or to secure the jack in a raised position in case the jack itself fails.\nTherefore, it would be desirable to have a lifting device that enables a user to lift the front or rear wheels of a vehicle. Further, it would be desirable to have a lifting device that provides additional safety to a user by holding a vehicle in a raised position even if the jack itself fails to hold. In addition, it would be desirable to have a lifting device that is compact and easy to store."}
{"text": "By virtue of such advantageous features as light weight, small thickness and low power consumption, liquid crystal display devices have been used in various fields as display devices of OA equipment, such as personal computers, and TVs. In recent years, liquid crystal display devices have also been used as display devices of portable terminal equipment such as mobile phones, car navigation apparatuses, game machines, etc.\nIn recent years, a liquid crystal display panel of a fringe field switching (FFS) mode or an in-plane switching (IPS) mode has been put to practical use. The liquid crystal display panel of the FFS mode or IPS mode is configured such that a liquid crystal layer is held between an array substrate, which includes a pixel electrode and a common electrode, and a counter-substrate, and switching is realized by rotating liquid crystal molecules of the liquid crystal layer in a plane parallel to the substrates. Such a display mode has an advantage of, for example, a wide viewing angle.\nIn the structure of the FFS mode or IPS mode, there has been a demand for an improvement of various display defects due to a displacement of an alignment direction of liquid crystal molecules from a desired direction. For example, in some cases, an image persistence phenomenon occurs on the liquid crystal display device of the FFS mode. The image persistence phenomenon is such a phenomenon that, for example, if an intermediate gradation image is displayed on the entire screen after a black-and-white checkered pattern is displayed on the screen, the checkered pattern slightly remains like an after-image. In recent years, various methods have been proposed for relaxing such an image persistence phenomenon."}
{"text": "1. Field of the Invention\nThe invention relates to anchoring of floating vessels. More particularly, it relates to a retrievable anchor for a floating offshore structure used for drilling or producing hydroacarbons located beneath a body of water.\n2. Prior Art\nThe following discussion of the prior art for illustrative purposes is centered on anchors used with tension leg platforms, Nevertheless, this invention is useful for other floating structures or vessels requiring anchors.\nA tension leg platform is sufficiently buoyant so that while floating on the water surface it keeps cables, respectively connected to the platforms and a fixity at the subsea bottom, under tension. One way proposed to provide the fixity adjacent the subsea bottom is through the use of deadweight anchors. Such anchors may include an enclosure with a concrete covering of predetermined thickness and weight; the thickness is a function of the initial ballast required to lower the anchor member to the ocean bottom.\nWith the anchor on the subsea bottom, ballast material, which may be either a flowable granular material (such as taconite or granular hematite) or sea water, is added to it. But, there are problems associated with each. The problem with granular material is that it is irretrievable and must be wasted at the drilling site. In the case of sea water, the problem is getting a large enough container that is still retrievable to hold sufficient water to provide restraint or resistance to movement as in a sea anchor.\nAnother way to provide fixity for the cables is to use solid steel or concrete anchors. They, however, upon relocation are either abandoned or lifted up by costly winch equipment.\nNeedless to say, wasting ballast such as the granular material or a concrete anchor is costly. These costs are eliminated by utilizing the present invention summarized below because it is easily retrievable. Consequently, it is possible to write off the fabrication and installation cost over a period involving several subsequent uses rather than just one use.\nFurther, the use of the present invention with small boats such as rowboats and sailboats should not be overlooked. The reason is that retrival of it does not require great physical strength of the retriever (for example a person in the boat), since the chain is first raised to the boat and then the empty shell. Other advantages of the invention will become obvious from the Summary of the Invention, Description of the Preferred Embodiment, the Claims and the drawings."}
{"text": "The field of the invention is in the art of firearms and more particularly in the art of adapters to provide for the firing of conventional low power ammunition in a gun originally designed for higher power cartridges.\nThe economic and psychological advantages of training, practicing, and the taking of small game with low power ammunition in a gun conventionally designed for higher power ammunition is well recognized. Adapter \"cartridges\", sometimes called auxilliary cartridges, for containing the low power round and fitting the chamber of the high power gun have been known since before the turn of the century. A conversion adapter for the Government model .45 caliber ACP converting it to fire .22 caliber rimfire cartridges has been available since prior to World War II. That conversion required a new barrel since the calibers of the bullets were different. The invention disclosed herein does not require any change of barrels since the bullet diameters of the original high power cartridges and the conversion low power cartridge are substantially the same.\nTypical examples of modern prior art devices are exemplified by U.S. Pat. No. 3,771,415 to patentees Into and Costello, and U.S. Pat. No. 3,776,095 to patentee Atchisson. A review entitled, \"AR-15 Rimfire Conversion\", by the Technical Staff of the National Rifle Association appearing in the American Rifleman for May 1973 commencing at page 63, is pertinent and informative of problems in the prior art devices."}
{"text": "Technical Field\nThe present disclosure relates to a colored contact lens. More particularly, the present disclosure relates to a colored contact lens with three-layered structure.\nDescription of Related Art\nNowadays, concerning fashion or clothes matching, consumers begin to wear colored contact lenses capable of altering their natural iris color. However these different colored or patterned contact lenses would fade or discolor so as to make the wearers have ocular allergy.\nA conventional colored contact lens is manufactured by forming a colored layer on the contact lens surface. However, when the consumer's finger touches the colored layer side during the process of wearing or cleaning contact lenses, the colored layer would easily peel from the contact lens surface. Therefore, the colorant in the colored layer is likely to fall off and to adhere on the surface of the wearer's eye so that the ocular allergy or infection of the wearer would occur.\nTherefore, a colored contact lens having an excellent colored performance is highly demanded."}
{"text": "Although induction motor drives still are the workhorse of industry, the switched reluctance motor drive has been actively researched over the past decade with very promising results. The switched reluctance machine has a simple and rugged construction as well as very good overall performance over a wide torque-speed range. Recently, doubly-salient switched reluctance motors have been found to be an attractive alternative to more conventional synchronous and induction machines in low horsepower converter fed variable speed drive applications.\nBecause the current waveform of the switched reluctance motor must be carefully programmed to extract the maximum torque per ampere, this machine is more accurately termed a current regulated stepping motor (CRSM or CRS motor). The fundamental feature of this type of motor drive is that the CRS motor requires only a unidirectional current and thus the drive circuit topology and corresponding switching algorithm is greatly simplified. A detailed comparison of the CRS motor drive with a high efficiency induction motor drive has indicated that performance parameters such as torque per unit stator volume, torque per unit inertia and torque per unit copper weight, can be made equal to that of an induction machine, or in some cases, even exceed the induction machine. M. R. Harris et al., \"A review of the Integral Horsepower Switched Reluctance Drive,\" IEEE Trans. Ind. Appl. Soc., Vol. IA-22, July/Aug. 1986, pp. 716-721.\nWhile the recent work on CRS machines has yielded encouraging results, in several respects the machine is less than optimum. For example, taking an eight stator/six rotor pole CRS motor, with typical excitation of one of the four phases, only one quarter of the stator inner circumference is utilized to make a contribution to torque development at any instant. Secondly, the inductance variation of the occupying coils over each one fourth of the machine is limited by the so called double salient design. To have a comparable power rating to that of the conventional induction machine of the same size with such a limited air gap inner surface area, the CRS machine is designed to operate in a deeply saturated condition. The corresponding active material is thus under severe electromagnetic stress. It is apparent that any method of ensuring that the other three quarters of the inner circumference of the stator remain \"active\" would be a very significant step toward improving torque production and therefore, towards relaxing the severe electromagnetic stress and consequent iron losses in the active materials of the machine.\nElectromechancial energy conversion in a CRS motor is accomplished by means of a time varying inductance due to the temporal variation of the rotor position. The switched reluctance motor typically has phase winding coils located on opposite sides of the stator and a multi-pole rotor. The poles of the rotor are brought into and out of alignment with the poles of the stator as the rotor rotates, periodically increasing and decreasing the inductance of the coils for each phase. In an idealized machine, the inductance varies with rotor angle as a form of triangular wave, with current being supplied by the converter to the phase coil during the time that inductance is increasing to obtain motoring operation. Generating torque is produced if the machine is excited during the interval in which the inductance of the winding is decreasing. Torque production is proportional to the square of the current and therefore independent of current polarity if mutual inductance is not involved. Hence, the windings can be excited with unidirectional currents.\nIn general, to maximize the torque for a given current in a CRS machine, the change in inductance as a function of rotor position should be maximized. However, in a conventional CRS motor, the ability to maximize the change of inductance with rotor position is limited. To maximize inductance, the stator pole should be completely aligned with the rotor pole but to minimize inductance the stator pole should be totally nonaligned with the rotor pole. In general, this requirement, and limitations imposed by leakage flux and the permeance of the main flux, result in a maximum achievable ratio of maximum inductance to minimum inductance of ten to one, when saturation is neglected. Another limitation, set by the small working area of the double salient structure, is that the permeability of cross section is strongly affected by the excitation level, which is detrimental to maximizing the ratio of maximum to minimum inductance. Thus, a significant improvement in the change in inductance with respect to rotor position in a CRS machine is not practical without increasing the number of poles, which unfortunately reduces the power output or speed of the machine or, conversely, requires very high switching frequencies in the associated power converter."}
{"text": "In the manufacture and production of processed meat food products, the final amount of ground meat components in the final food product is labeled, such as the percent fat and/or amount of protein, for example. The amount of ground meat added to the process is used to calculate the amounts of other additional ingredients that are based on the amount of meat added. Therefore, there is a need to measure the amount of ground meat transported into the process hopper from the meat hopper and added to other ingredients. One way to measure the amount of ground meat is to monitor the flow rate of the system. However, flow rates of non-liquid streams, such as ground meat, can be difficult to measure in a continuous process.\nOne method of measure utilizes a proportional-integral-derivative controller, or PID loop control, often used as a feedback loop control in process manufacturing systems. The PID controller obtains a measured value from the process and compares it with a setpoint or target value. The difference between the two values, if any, is an error value, which is then used to calculate a correction to the process input variable so that the correction will remove the error from the process output measurement and bring it closer to the target setpoint.\nIn a ground meat processing system, loss in weight data is measured to provide PID loop control to establish the flow rate of ground meat. Loss in weight data is obtained by measuring the weight of the ground meat stream fed to the system (i.e., fed to the pump) and subtracting the final weight of the stream at the same location after a set time period has elapsed. When considered in conjunction with the elapsed time, this yields an actual flow rate which can be used to obtain a difference in flow rate between the actual flow rate value and the target value, which is an error value. The pump speed is then adjusted based on the loss in weight data provided to the PID loop control.\nOne problem with the PID loop control being based on loss in weight data, however, is that as the hopper is being filled with the meat, the loss in weight data cannot be properly calculated, as can be seen in FIG. 2. This is because meat is being added to the hopper at the same time that meat is leaving the hopper and being transported through the pump, therefore, an accurate measurement of weight loss cannot be properly measured at the same time. As a result, during times of filling the hopper flow rate measurements must be temporarily halted due to the rapid increase in weight within the hopper. Therefore, the halt in collecting loss in weight data interferes with the operation of the PID loop control and during these periods of filling there is no PID loop control.\nWithout accurate and continual PID loop control, or any other type of flow rate control of the system, inaccurate flows of the meat stream may cause the amount or weight of meat added to the final meat product to be inaccurately calculated."}
{"text": "Light emitting diode devices are improved remarkably in terms of light emitting efficiency and are free from mercury which is a toxic substance, whereby use thereof for illumination instead of conventional fluorescent lamps has energetically been studied in recent years. It is preferred that an illumination source emits a light having a spectral distribution close to that of sunlight. Thus, it is preferred that the illumination source has a flat spectral distribution containing components with every wavelength in the optical region. As white light source using light emitting diode devices, those obtained by a combination of three light emitting diodes for red, green, and blue emission or a combination of a light emitting diode for blue to ultraviolet emission and a photoluminescent material which absorbs the light emission, to thereby emit light from red to green, have been used. Examples of the former include so-called a three-in-one LED light source in which LEDs for red (AlInGaP), green (InGaN) and blue (InGaN) are included in one unit, and examples of the latter include a light source in which an InGaN-based blue LED and a so-called YAG fluorescent material to absorb a blue light to thereby produce a yellow fluorescence are combined."}
{"text": "The present invention relates to medical apparatus and more particularly to electroconvulsive therapy apparatus."}
{"text": "As a method of transmitting a large volume of data at a high speed, a system combining an OFDM (Orthogonal Frequency Division Multiple) and CDMA is under study in recent years. As a system combining OFDM and CDMA, there are two schemes; a scheme whereby chips with spread data are mapped on subcarriers in a frequency domain and a scheme whereby chips with spread data are mapped on subcarriers in a time domain.\nWhen data is spread in the frequency domain, there may be violent propagation channel variation in the frequency domain caused by frequency selective fading due to a multipath environment, which causes orthogonality among spreading codes to be lost deteriorating its reception characteristic though a frequency diversity effect can be obtained during despreading.\nWhen data is spread in the time axis domain, a variation in the propagation channel in the time axis domain is relatively moderate compared to that in the frequency domain, and therefore little frequency diversity effect is obtained, but orthogonality among spreading codes is maintained. However, data assigned to subcarriers with sharp drops has a very low reception SNR, and therefore there is a high probability that the data may be completely erroneous.\nEspecially when codes are multiplexed using M-ary modulation such as 16QAM, deterioration of reception performance due to loss of orthogonality among spreading codes is drastic, and therefore spreading in the time axis domain has a better characteristic than spreading in the frequency domain."}
{"text": "Structural modifications for specific intended uses during manufacture of the workpieces are commonly required. For example, angle irons or structural angles (i.e., angle irons) often require that multiple holes be provided at various locations along the angles. Given the strength of the materials of the angles, significant forces are commonly required to create those holes, and thus significant mechanisms are required to generate those forces. Such forces are similarly required for other operations which might be required in which holes completely through the angle may not be required, for example, stamping identifying characters in the surface of the angle.\nMachines have, of course, been used which can create holes (e.g., by punching or drilling), or stamp identifying information, in such workpieces, usually in a facility where the workpieces are being worked on (e.g., where a long blank is being punched to provide whatever holes are required for the intended use of the part pieces, with individual elements being sheared from the blank to form the individual part pieces).\nWorkpieces such as structural angles which are not simply flat and/or are made of strong material can be particularly difficult to work with in creating holes. For example, structural angles may have two longitudinal members or legs connected at right angles along an edge (often by bending a single flat longitudinal member along a line extending in the longitudinal direction), and typically are made of strong metals such as steel or iron to provide the strength required in many construction and manufacturing applications. In order to create holes in both of the legs of structural angles, separate punches have been used for each the two different legs of the angle, with one punch for one leg of the angle and a separate punch for the other leg of the angle. Those punches have shared a mechanism which serves to properly position the angle lengthwise for punching (e.g., along the X-axis), and have their own separate drives to move each individual punch head assembly to the correct location (along the Y- and Z-axes) for punching a hole at the selected longitudinal location of the structural angle. Not only can the cost of such dual punches be significant, but the speed of operation is also impacted since clearance requires that the punch head assemblies be spaced along the X-axis, resulting in time being required to move the entire structural angle along the X-axis for punching holes in both legs of the angle, even if the holes are at the same position along that X-axis. Further, precise positioning of holes which are supposed to be at the same longitudinal position on the angle may not be achieved if the structural angle is not moved accurately along its X-axis between the different punches.\nAdditionally, the operation of punch mechanisms for structural angles used heretofore have either resulted in inefficient handling of the structural angles or required additional mechanisms for the punch. That is, some machines have positioned the die portion of the punch so that it is always positioned so that it will be contacting one side of the angle (as desired during punching operation). However, such positioning necessarily results in the angle dragging on the punch die when it is moved through the mechanism. While such undesirable dragging has been avoided by either using support rollers which swing up to support the material clear of the die when moving the angle, or by providing a separate mechanism (e.g., separate hydraulic cylinders) to lift the punch die against the angle when punching is required, such lift methods have required separate items (hydraulic cylinders, rollers, tapered slide blocks, hydraulic valves, and the associated IO and timing) to make things work correctly.\nThe present invention is directed toward overcoming one or more of the problems set forth above."}
{"text": "Styrenic polymers such as ABS are known to be immiscible with polyamides, particularly in the blend ratios of interest ranging from 1:9 to 9:1 (see e.g. Y. V. Lebedev, et al, J. Applied Polymer Science 25, 2493, 1980).\nDue to this mutual incompatibility, simple mechanical blends of ABS and nylon 6 tend to be brittle and readily delaminate when the molded parts are flexed or twisted. U.S. Pat. No. 4,496,690 discloses compositions of modified ABS and nylon blends in which the ABS portion was significantly modified through copolymerization with as much as 10% acrylamide monomer units to obtain useful levels of compatibility and properties.\nEuropean patent application #202214 (1986) teaches the use of styrene-acrylonitrile-maleic anhydride terpolymers as compatibilizers for ABS-nylon blends."}
{"text": "In a process of manufacturing an optical element, such as a light guide plate, a diffusion plate and the like, a vacuum-lifting device is used for vacuum-lifting the optical element to a preset location to be machined."}
{"text": "It is known to warp knit a support tape for a slide-fastener stringer half. This tape has a pair of parallel longitudinal edges and a pair of opposite faces. A monofilamentary coupling element having a succession of like turns is laid on one face at one edge thereof and stitched to the tape. The stitching is typically a multithread double chain stitch or double-lock stitch of class 400 (Federal Standard 751a).\nSuch stringer halves have shown themselves to give long service. They have relatively good transverse and longitudinal dimensional stability and are relatively wear resistant. Nonetheless in most of these arrangements the warp yarns, that is the filaments running mainly longitudinally, form chains into which a group of weft yarns, that is filaments running mainly transversely, are laid. These weft filaments are invariably laid in over at least four chains (see German open application DT-OS 2,016,141 and U.S. Pat. No. 3,708,836), so as to give maximum transverse dimensional stability.\nMore particularly, in this arrangement, the wales are of the single-bar type formed as warp chains which are bridged only by the filler or weft. As indicated the sole connecting threads or yarns between the warp chains are the laid-in weft filaments or yarns which extend over four warp chains each so that four such weft inlays fill each warp loop to provide the transverse stability of the tape (stability in the plans of the tape perpendicular to the warp chains) and, in addition, prevent longitudinal extension and contraction (change of dimension in the plane of the tape parallel to the warp chains). The dimensional stability is most pronounced if the tapes are composed of synthetic-resin yarn. which has been thermofixed, i.e. subjected to a heating or ultrasonic treatment which relaxes internal stresses within the yarn.\nThe advantage of single-bar warp chains, i.e. individual spaced apart wales of loops interconnected only by the weft, is that channels are formed between wales so that the row of stitching for attaching the coupling element can lay along one of these channels or valleys.\nWhile such an arrangement has proven itself relatively desirable it has been found to be advantageous to operate the sewing machine which secures the coupling element to the tape edge at still higher rates of speed. Furthermore, such a knit does not fully utilize the capacities of the industrially normally used warp-knit machine, so that capacity is lost."}
{"text": "There is a growing market and need for electric vehicles that can be fueled using renewable energies such as solar, wind, hydro and biochemical (fuel cell) processes that produce electricity. While there are existing renewable energy systems for supplying electricity to homes and businesses, there are no currently no known systems adapted for use as personal vehicle fueling and storage stations.\nCurrently, the electric vehicle market is targeting the use of hydrogen fuel cells for electric refueling in accordance with renewable energy priorities recommended by the U.S. Government. Larger scale hydrogen fueling stations have been proposed in government demonstration projects in preparation for building a national infrastructure. In January 2001, HyGen Industries LLC and Hamilton Sundstrand Space Systems International, Inc., in Torrance, Calif., undertook a project to develop, demonstrate, deploy, commercialize and market renewable hydrogen generation for vehicle fueling systems. In July 2001 the first hydrogen production and fueling station in the Los Angeles area started operation at Honda Motor Company's research and development center. In October 2001 a hydrogen refueling station was installed in Brunnthal-Nord, near Munich, for refueling of a small test fleet of fuel-cell powered buses to connect south-east localities with the east railway station of Munich. The technology was provided by ET (Energy Technology), a company founded in 1997 by hydrogen experts from the space company DASA. The Munich installation was the fourth in Germany after previous installations at the Munich Airport and in Hamburg and Oberstdorf. In 2003, the Schatz Hydrogen Generation Center installed a complete hydrogen production, compression, storage, and dispensing facility at SunLine Transit Agency in Thousand Palms, Calif. The facility provides compressed hydrogen at 3600 psig for use in the SERC hydrogen fuel cell-powered vehicle fleet operated by SunLine Transit and the City of Palm Desert.\nHowever, such current efforts to create renewable energy refueling stations face a number of intractable problems. Stations which take in compressed hydrogen gas for generation of electricity for refueling electric vehicles have problems with the cost and safety of transporting remotely-generated hydrogen to the station, and the large infrastructure costs of providing sufficient on-site electrical energy storage for use in recharging electric vehicles. Vehicles that come to a central refueling station may experience a long refueling delay to recharge vehicle batteries or to take in hydrogen gas for on-board fuel cell generators. Long delays in refueling each vehicle multiply into worse problems of long lines and traffic congestion for vehicle operators. It is deemed highly desirable to provide a renewable energy refueling approach using facilities that are distributed and adapted to refuel electric vehicles on a personal scale."}
{"text": "1. Field of the Invention\nThe present invention relates to photosensitive compounds, more particularly, derivatives of halomethyl-1,3,5-triazines.\n2. Discussion of the Prior Art\nCompounds that decompose to generate free radicals (free radical generating agents) upon exposure to light are well known in the graphic arts. Organic halogen compounds, which are capable of generating free radicals such as a chlorine free radical or a bromine free radical upon exposure to light, have been widely used as photoinitiators in photopolymerizable compositions, as photoactivators in free radical photographic compositions, and as photoinitiators for reactions catalyzed by acids formed by light. The spectral sensitivity of these compositions may be broadened by the addition of sensitizers, which, in essence, transfer their absorbed energy to the organic halogen compound. The use of such halogen compounds in photopolymerization processes and free radical photographic processes have been described in Kosar, Light-Sensitive Systems, J. Wiley & Sons (New York, 1965), pp. 180-181, 361-370.\nHalomethyl-1,3,5-triazines are known to be initiators for a number of photochemical reactions. They are employed to produce free radicals for initiating polymerization or color changes and for initiating secondary reactions upon liberation of acid by the interaction of the free-radicals when hydrogen donors are present.\nExamples of the use of halomethyl-1,3,5-triazines in the free radical polymerization of acrylate monomers are described in U.S. Pat. No. 3,905,815; U.S. Pat. No. 3,617,288; U.S. Pat. No. 4,181,752; U.S. Pat. No. 4,391,687; U.S. Pat. No. 4,476,215; and DE 3,517,440. U.S. Pat. No. 3,779,778 discloses the photoinitiated acid catalyzed decomposition of pyranyl ether derivatives to produce photosolubilizable compositions useful as positive printing plates. Chromophore substituted styryl-1,3,5-triazines and their uses are disclosed in U.S. Pat. No. 3,987,037 and U.S. Pat. No. 3,394,475.\nRadiation sensitive compositions containing bi- and polyaromatic substituted triazines are disclosed in U.S. Pat. No. 4,189,323."}
{"text": "The present invention is related generally to fiberoptic cables and waveguides, and optical coupling therebetween. More specifically, the present invention relates to a method for making the endfaces of a plurality of fiberoptic waveguides on a substrate have a smooth surface for substantially reducing light scattering, and providing high-efficiency coupling with the end of a fiberoptic cable or optical fiber.\nIt is known in the art to use semiconductor substrates for providing a support for an optical waveguide formed on a substrate. Typically, one end of the waveguide is optically coupled to an end of a fiberoptic cable or optical fiber for permitting light waves to travel therebetween. Coherent light waves such as produced by lasers, are employed in fiberoptic communication systems, for example, for transmitting data signals through use of optical waveguides to couple the light waves to and from fiberoptical cables to photodetector and other circuitry.\nA major problem providing high-efficiency optical coupling between optical waveguides and fiberoptic cables or optical fibers. It is important to achieve as near perfect alignment near the end of the fiberoptic cable and the endface of the waveguide for maximizing the transfer of light therebetween. Also, it is important to avoid any surface irregularity on the endface of the waveguides that may cause light scattering.\nIt is known in the prior art to form one or more V-grooves in  less than 100 greater than  silicon wafers, for retaining an optical fiber or fiberoptic cable, for coupling to the endfaces of waveguides cores formed in an opposing portion of the substrate. Note that although silicon (Si) is typically used for the substrate, other substrate materials include but are not limited to Gallium Arsenide (GaAs), and Indium Phosphide (InP).\nIn forming optical waveguides on various substrates, it is known in the art to employ different fabrication processes depending on the substrate material. The fabrication process may incorporate any one of chemical vapor deposition (CVD), fire flame deposition, and so forth. In such fabrication, V-grooves may be formed via chemical etching on one portion of the substrate. In this regard, anisotropic etchants for silicon substrates include potassium hydroxide (KOH), KOH and alcohol, ethylene diamine, and so forth. Etchant solutions for providing an anisotropic etching of GaAs or InP substrates include solutions of sulphuric acid (H2SO4), hydrochloric acid (HCl), bromic acid (HBr), and so forth. During such processing, masking is required for producing the proper formations, whereby the mask must have a relatively low etch rate. Examples of methods for forming V-grooves and optical waveguides in different portions of semiconductor substrates, and methods for the use of such grooves in retaining a portion of an optical fiber or cable, and aligning the same with an associated waveguide having an endface opposite the groove are shown in U.S. Patents of Carney U.S. Pat. No. 4,466,696, Welbourn U.S. Pat. No. 5,342,478, Menigaux et al. U.S. Pat. No. 5,518,965, Kitamura et al. U.S. Pat. No. 5,548,673, Yamane et al. U.S. Pat. No. 5,557,695, and Mizuta et al. U.S. Pat. No. 5,961,869. The teachings of these references are incorporated herein to the extent that they do not conflict herewith.\nMuch effort has been made in the prior art to develop better methods for coupling optical fibers or fiberoptic cables to waveguides formed on a semiconductor substrate. Examples of methods and apparatus for accomplishing such coupling are provided in many U.S. Patents, including Boivin U.S. Pat. No. 3,774,987, Blonder U.S. Pat. No. 4,966,433, Stein U.S. Pat. No. 5,046,809, Bossler U.S. Pat. No. 5,357,593, Bresby U.S. Pat. No. 5,371,818, and Harpin et al. U.S. Pat. No. 5,787,214. The teachings of the previously identified U.S. Patents are incorporated herein by reference to the extent that they do conflict herewith.\nThe teachings of Menigaux U.S. Pat. No. 5,518,965, entitled xe2x80x9cProcess For Producing A Structure Integrating A Cleaved Optical Guide With An Optical Fiber Support For A Guide-Fibre Optical Coupling,xe2x80x9d issued on May 21, 1996. This reference teaches the formation of a V-groove in one end portion of a semiconductor substrate in one embodiment, and also in the opposing end portion in another embodiment. An optical waveguide is formed on the portion of the substrate not containing a V-groove. A method is taught for cleaving the ends of the waveguide for forming the endface or endfaces thereof, for coupling to the ends of optical fibers.\nIt is known in the art to provide an integrated optic chip with one or more V-grooves formed in one portion of a substrate, and one or more waveguides formed in another portion of the substrate, with endfaces of the waveguide cores opposing respective ones of the V-grooves. An example of a substrate formed with a plurality of V-grooves and waveguide cores, as indicated, is shown in FIG. 1. More specifically, as shown, a substrate 1 has a plurality of juxtaposed V-grooves formed in a top first portion of the substrate 1. In an opposing top second portion of the substrate 1, a cladding layer 5. A plurality of waveguide cores 7 are formed within the cladding, as shown, with their longitudinal axes parallel to one another and to the parallel axes of the V-grooves 3. In the example shown, a dicing saw is used to form a transverse rectangular open slot 9 between the portions of the substrate 1 containing the V-grooves 3 and the waveguide cores 7. In this manner, the sloping endface walls of the V-grooves 3 are removed to provide vertical walls, for insuring that the ends of optical fibers 11 carried in the V-grooves 3 can abut endfaces 13 of the waveguide cores 7. However, such use of a dicing saw produces a rough surface on the endfaces of the waveguides 7, making it difficult to obtain high-efficiency coupling between the waveguide cores and optical fibers 11 carried in the V-grooves 3. Good coupling is precluded by light scattering caused by the rough surface on the waveguide endfaces 13 and sidewall 15 of the substrate 1. Polishing of the endfaces 13 of the waveguide cores 7 can serve to reduce the light scattering, and improve the efficiency of optical coupling thereto. However, such polishing is extremely difficult.\nAn object of the invention is to provide an efficient method for forming a smooth endface for a plurality of waveguides formed on a semiconductor substrate along with opposing V-grooves.\nAnother object of the invention is to provide a method for simultaneously cleaving a portion of the coupling ends of a plurality of waveguide layers or cores formed on a semiconductor substrate, in a manner avoiding cleaving of the substrate.\nYet another object of the invention is to provide a method for efficiently and simultaneously cleaving coupling end portions of a plurality of waveguide cores formed in cladding on a semiconductor substrate, for providing a smooth endface on the waveguide cores, thereby providing high-efficiency coupling to the ends of associated optical fibers or fiberoptic cables, respectively.\nWith the problems of the prior art in mind, the above-identified objects and other objects of the invention are satisfied in one embodiment by a method including the steps of initially forming a plurality of V-grooves in one top portion of a semiconductor substrate, with a plurality of waveguide cores equal in number to the number of V-grooves formed within a cladding layer or layers on another top portion of the semiconductor substrate, with endfaces of the waveguide cores opposing respective ones of ends of the V-grooves. A dicing saw or other means is used to form an open slotway between the portion of the substrate V-groove or grooves, and the portion of the substrate with the waveguide cores. In one embodiment, after using a dicing saw to form the open slotway, the portion of the substrate with the V-grooves is masked along with a portion of the open slotway. Then, an etchant is used to undercut through the sidewall of the portion of the substrate immediately below the cladding and waveguide core layers for undercutting the latter to provide a cantilevered waveguide portion. The mask material is removed, and a scribe line is formed across the transverse portion of the top surface of the cladding and waveguide core layers, with the scribe line preferably being parallel to the edge of the undercut portion of the substrate immediately underlying the cladding and waveguide core layers. Next, a downward force directed toward the V-grooves is applied to the cantilevered portion of the waveguide forward of the scribe line. The force cleaves the cladding and waveguide core layers, providing a smooth endface for the waveguide cores opposing respective V-grooves.\nIn a second embodiment of the invention, the open slotway is cut into the substrate at an angle via a dicing saw, followed by the aforesaid masking, etching, scribing, and cleaving steps.\nIn a third embodiment of the invention, the scribe line is formed on the top cladding layer forward of the top edge of the underlying substrate proximate the undercut portion of the substrate. In this third embodiment, a force is applied to the cantilevered portion of the cladding and waveguide core layers, causing these layers to breakaway leaving the endfaces Of the cladding and waveguide core layers at an angle.\nIn a fourth embodiment of the invention, the angled slotway is formed as previously indicated, but the scribe line is formed on the top cladding layer in alignment with the top edge of the underlying substrate proximate the undercut portion. This step is followed by masking, and then etching for forming the undercut. Next, the scribe line is formed in alignment with the top edge of the underlying substrate proximate the undercut portion, followed by cleaving to leave a smooth and vertically oriented endface for the cladding and waveguide core layers.\nIn a fifth embodiment of the invention, the cleaving force is provided by gluing the end of an object to the end of the cantilevered portion of the cladding and waveguide core layers. Setting of the glue is followed by pulling the object to apply a tension force on the cantilevered end (the force is longitudinally directed away from this end). Simultaneously a downward directed force can also be applied to the object proximate the region where it is connected to the cantilevered end of the waveguide and cladding layers, for cleaving the latter.\nIn a seventh embodiment of the invention, a high stress layer is formed on top of the cladding, with the front edge of the high stress layer being aligned with the edge of the substrate proximate the undercut region thereof. The high stress layer can be formed from an appropriate metal, such as nickel. Thereafter, the cantilevered portion of the cladding and waveguide layers can be cleaved by applying a downward force. In the latter embodiment, it is preferred but not absolutely necessary to use a scribe line before cleaving."}
{"text": "As an extremely complex system, petroleum comprises not only alkanes, cyclic alkanes and aromatic hydrocarbons of different molecular weight, but also a small amount of sulfides, oxides and nitrides, as well as some trace amounts of metal compounds of iron, copper, nickel, vanadium, etc. In the process of petroleum processing, these hydrocarbon or non-hydrocarbon substances will subject such chemical reactions as cracking, condensing and coking. Now, the catalytic hydrogenation process is usually conducted by adding hydrogen into the feedstock oil under conditions of high temperature and high pressure both at home and abroad. However, with the ever-increasing contradictions between the trend of heavy crude oil and need to be light for the petroleum products, the feedstock oil to the reprocessing device is becoming heavier and inferior in quality, that finally causes some serious problems to the catalyst for hydro cracking, hydro desulfurization and catalytic cracking devices in many refineries and petrochemical plants such as coking, scaling and eventually deactivated, even influencing the catalytic bed of hydro refining and hydro cracking to fasten pressure decline, thus seriously affecting and confining the action cycle and stable operation of the devices.\nThe problem of catalyst deactivation is obvious especially in the hydro cracking device because of its high reacti n pressure of 15-20 MPa, high temperature of 380-420° C. and the one-lot filling of catalyst of 500-1000 tons. In some cases, the equipment has to be stopped to maintain only after 3 to 5 months' operation due to catalyst deactivation. Thus, it is an urgent task to prevent or alleviate the catalyst deactivation and enhance catalyst efficiency in hydro cracker to lengthen the action cycle to acquire more economic benefits.\nThe catalysts used in the process of petroleum hydrogenation are inorganic, while the feedstock involved in the hydrogenation reaction are organic materials and non-polar hydrogen. Accordingly, the contacting activity of feedstock with catalyst is relatively poor, and coke easily formed on the surface of the catalyst will affect activity of the catalyst.\nThere have been quite a few reports both at home and abroad on study of new hydrogenation catalysts in recent years. However, none of them addresses to the addition of catalytic active accelerant to the feedstock oil during the hydrogenation to improve the catalyst activity, the depth of catalytic hydrogenation and the yield of catalytic hydrogenation products, and weaken the reaction conditions of catalytic hydrogenation. It would be a tremendous contribution to the international catalytic hydrogenation techniques that catalytic active accelerant is applied in practical production to increase the catalyst activity and prolong the service life of the catalyst."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for the preparation of zero standard sera useful in immunoassay techniques.\n2. Description of the Prior Art\nRadioligand hormone assays have now become routine in hospital laboratories throughout the world. The key to the sensitivity of the test is the availability of hormone-free sera to be used as controls calibrators and standards and also, in some cases, as diluents.\nFor example, standard controlled serum which contains no thyroxine is necessary in order to obtain the zero point for the standard curve used in radioimmunoassays or immunoradiometric assays or enzyme immunoassays. In addition, the use of zero standard control serum is also necessary to obtain accurate readings of hyperthyroid serum. Thus, when zero standard control serum is available, hyperthyroid serum may be diluted by a known factor with the zero standard control serum. T4 tests may then be run by the above-identified procedures and compared with values obtained from standard curves. After the readings are obtained from these curves, dilution factors are applied to the readings. Accurate factors are only obtained when the zero standards have very low levels of hormones.\nEisentraut, U.S. Pat. No. 3,776,698, for example, prepares zero standard serum by contacting a quantity of acidified serum with a particulate inorganic crystalline sorbent material which is selected from the group consisting of phosphates, oxides, hydroxides, silicates, carbonates, aluminates, and sulfates, of the metallic elements in groups IA, IIA, IIIA, IIB and VIII of the periodic table.\nTwo alternative methods of preparing hormone-free serum are also presently used: 1) TSH free serum can be made by accumulation of a pool of serum from patients with hypopituitarism who are on replacement thyroxine therapy, or 2) serum can be made by extended high speed centrifugation of free carbon-adsorbed serum from a healthy population.\nThe use of free carbon, however, carries with it difficulties in time and equipment, when large volumes of serum are to be processed.\nIn order to be able to effectively remove thyroid hormones (such as T3 or T4) from the serum, it may be necessary to provide increased surface area in the treating carbon. Commercially available carbon is usually made by activating the same with steam at temperatures around 1000.degree. C. The activation probably removes tars and insertion materials trapped between the naturally layered structure of the carbon, thus allowing access to the interlayer regions. The gain in surface area which is obtained, however, results in sacrifice of the structural rigidity of the final product.\nThe mechanical forces imposed on mixing the carbon with serum at high speed centrifugation are sufficient to fragment the fragile activated carbon structure. The supernatant serum normally contains carbon fragments visible as the appearance of dark color. Removal of the tiny carbon fragments from the serum requires the addition of technical grade kaolin, which acts as a gathering agent, followed by further mixing and centrifugation. Sometimes, even colloidal flocculation is not satisfactory for complete removal, and filtration through a micron size membrane is required.\nLarge carbon granules (having a range of about 400 microns in diameter) show poor efficiency of thyroid hormone removal, whereas smaller powdered carbon granules used in batch, while more efficient, are inconvenient given the slowness and extra steps involved in separation procedures.\nAn attempt to increase the efficiency of hormone removal, by artificially providing a pressure differential across a bed of carbon particles, runs into various problems. The structural rigidity of the carbon particles may be insufficient to withstand the applied pressure, and the layered arrangement may break, with concommitant collapse of the particles. Also, channeling quickly occurs across the bed, followed by rapid loss in removal efficiency.\nCarbon has been loaded on a sheet having particles homogeneously distributed and firmly retained therein, and having substantial freedom from ruboff, flaking and smudging, as described in Bodendorf et al, U.S. Pat. No. 3,149,023. The Bodendorf et al sheets, however, are used as cigarette filters, air filters, gas filters, wrappers for fruit and substances prone to discoloration or spoilage by gases in the atmosphere, deodorizer layers in laminated sheet products for sanitary napkins, and for surgical dressings for wounds, and the like. They are not used in the removal of thyroid hormones from serum.\nThe Bodendorf materials are not even described as being useful for increased pressure applications and, most importantly, for increased pressure applications in the treatment of a biological material such as serum, wherein the integrity and composition of the final product is critical.\nIn preparing zero standard serum, the final product should contain all the necessary natural serum components except the hormones. Inasmuch as possible, proteins, carbohydrates, lipids, vitamins, and salts should remain at substantially the same levels as prior to the treatment. There is no indication in Bodendorf et al that their carbon filled sheets can be used for the treatment of such delicate biological materials as human serum, that the treatments could be carried out at high pressures, and that they would result in selective removal of hormones.\nIn sum, a method for the removal of hormones from human serum to prepare zero standard serum which method is quick, inexpensive, and simple to operate, which uses materials which are simple to prepare, and which is non-destructive and non-modifying of the final standard product has not yet been described by the prior art.\nA need therefore continues to exist for such a method, and the present invention fulfills that need."}
{"text": "The present invention relates to a recorder for electric appliances. More particularly, the invention relates to a recorder for electric appliances for recording the time during a predetermined period of time that an electric appliance is operating.\nObjects of the invention are to provide a recorder of simple structure, which is inexpensive in manufacture, installed with facility and convenience, used with facility and convenience, and functions efficiently, effectively and reliably to record the time during a period of time that an electric applicance operates. This enables a study of the operation of the applicance so that it may be controlled for maximum efficiency and conservation of energy."}
{"text": "The present invention generally relates to the field of wireless communication systems and to systems and methods for managing wireless picocells.\nGrowth in wireless communication continues to increase. Demand for data services with high data bandwidth requirements has led to the introduction of multiple modulation techniques for wireless communication, such as Long Term Evolution (LTE), High-Speed Downlink Packet Access+(HSDPA+), and CDMA2000 1xEV-DO (Evolution-Data Optimized or “EVDO”).\nAdditionally, deployment of small cells including picocells and femtocells has become increasingly desirable for providing coverage. Small cells may be deployed, for example, in areas having high user density, such as airports or event venues. A small cell deployment typically has a 100 meter to 1 kilometer radius. Both voice and data modes are desired in small cell deployments. Development of multi-modal multi-modulation capable picocells is complex. Such picocell systems need management of backhaul and core network connectivity as well as mode-to-mode communications and synchronization of picocell features."}
{"text": "During filming sequences with a motion picture or video camera, it is often necessary for the camera to be moved over a floor or ground surface. This movement must be accomplished smoothly. Even small amounts of shock or vibration transmitted to the camera lens can result in shaky or erratic recorded images.\nOne technique for achieving smooth camera movement is to support the camera on a camera cart, dolly or crane which moves or rolls over rails or track. However, installing track is very time consuming. Placement of the track may also interfere with camera angles, since in most filming sequences, it is important that the track must remain out of the field of view of the camera. In addition, after the track is laid, the dolly is necessarily restricted to movement purely along the track. Even small changes to the path of movement of the camera (as may be spontaneously requested by the director, cinematographer, camera operator, etc.) cannot be achieved without changing the track layout. As a result, track is generally only used when absolutely necessary, such as where smooth movement over a very irregular surface (such as sand, gravel, unpaved outdoor surfaces, etc.) is needed.\nDue to the need for an extremely stable camera platform, even relatively smoother surfaces, such as paved surfaces, or indoor flooring, can create difficulties in achieving satisfactory filming. The motion picture or video camera is typically mounted on a camera crane or dolly having wheels which roll over the floor or supporting surface. Irregularities in the floor or supporting surface (for example, cracks, bumps, depressions, debris particles, etc.) generate shock and/or vibration impulses which can be transmitted through the structure of the dolly to the camera lens. These impulses degrade steady filming ability. In addition, since the camera is typically positioned substantially above the floor surface, the effects of a small crack or bump in the floor surface can be magnified at the camera position. To reduce the effects of irregularities in the floor surface, camera dollys and cranes having soft tires have been used. However, soft tires cannot eliminate the effects of larger irregularities in the floor surface. In addition, soft tires increase rolling friction, and make the dolly more difficult to move or steer.\nThe shock and vibration isolator described in Chapman, U.S. Pat. No. 4,989,823, incorporated herein by reference, has successfully been used for several years. This isolator design is very effective in isolating the camera from lateral (side to side) and longitudinal (front to back) shock and vibration impulses. It has been widely used with very good results on camera dollys. However, this isolator design has only limited ability to prevent vertical shock or vibration impulses from reaching the camera.\nAccordingly, there is a need for an improved device and system to allow for smooth filming sequences.\nMany camera cranes and dollies have a leveling system for keeping the camera platform in a level orientation as the crane or dolly arm pivots up or down. Typically, these leveling systems use a parallelogram linkage, with examples described in Chapman U.S. Pat. Nos. 4,360,187; 5,816,552; 6,345,919 and 6,450,706, incorporated herein by reference. These parallelogram leveling systems keep the camera platform locally level, that is level or parallel with the vehicle, such as a crane or dolly that they are mounted on (and not necessarily level with the ground). However, when the dolly rolls over an irregularity, such as a bump or crack, the dolly is momentarily inclined at a slight angle. This movement, or elevation angulation, causes the camera lens to also be momentarily inclined at a slight angle. This disturbance of the camera lens can become unacceptable, especially when filming distant subjects with telephoto lenses. Accordingly, there is a need for an improved camera platform leveling system able to better compensate for angulation."}
{"text": "Field of the Disclosure\nThe present disclosure, for example, relates to techniques for acquiring measurements of a shared spectrum and performing channel selection for access points using the shared spectrum.\nDescription of the Related Art\nWireless communication systems are widely deployed to provide various types of communication content such as voice, video, packet data, messaging, broadcast, and so on. These systems may be multiple-access systems capable of supporting communication with multiple users by sharing the available system resources (e.g., time, frequency, and power). Examples of such multiple-access systems include code-division multiple access (CDMA) systems, time-division multiple access (TDMA) systems, frequency-division multiple access (FDMA) systems, and orthogonal frequency-division multiple access (OFDMA) systems.\nBy way of example, a wireless multiple-access communication system may include a number of access points (e.g., bases stations, eNBs, and/or Wi-Fi access points), each simultaneously supporting communication for multiple UEs. An access point may communicate with UEs on downlink channels (e.g., for transmissions from an access point to a UE) and uplink channels (e.g., for transmissions from a UE to an access point).\nSome modes of communication may enable communications with a UE over different radio frequency spectrum (e.g., licensed spectrum and/or shared spectrum). With increasing data traffic in cellular networks, the offloading of at least some data traffic to a shared spectrum may provide a mobile network operator with opportunities for enhanced data transmission capacity. Prior to gaining access to and transmitting data over the shared spectrum, a transmitting apparatus may, in some examples, perform a listen before talk (LBT) procedure to gain access to the shared spectrum. An LBT procedure may include performing a clear channel assessment (CCA) to determine whether a channel of the shared spectrum is available. When it is determined that the channel of the shared spectrum is not available (e.g., because another device is already using the channel of the shared spectrum), a CCA may be performed for the channel again at a later time.\nBefore performing an LBT procedure to gain access to a channel of a shared spectrum, a transmitting device may need to select the channel from a plurality of channels for communicating over the shared spectrum."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital camera having a built-in printer in which a printing light is emitted to an instant film in accordance with image data, which is obtained by photographing a subject image, to perform hard copy of the image data.\n2. Description of the Related Art\nA digital camera utilizing a CCD image sensor and so forth is widely used. In the digital camera, an optical subject image is converted into an electronic image signal. Further, the electronic image signal is digitally converted for obtaining image data to be stored in a memory. Some of the digital cameras have a built-in printer to produce a photo print of the desired image data by means of hard copy. This image data is selected among the image data stored in the memory, and the hard copy thereof is performed on an instant film. The printer is provided with an optical head for emitting printing lights of respective colors of red(R), green(G) and blue(B). The optical head is disposed so as to confront a photosensitive surface of the instant film, and is moved on a parallel therewith. Upon selecting the desired image data for the hard copy and upon starting a print operation, the optical head performs an exposure in accordance with the image data, moving in a sub-scanning direction. The exposure is performed in a three-color frame-sequential manner to print a full-color image on the instant film.\nIn the above-described digital camera having the built-in printer, it is possible to easily perform the hard copy of the image data stored in the memory. However, even though the instant film on which the image is recorded by the hard copy exists, it is impossible to perform the hard copy of the same image again when the image data is erased from the memory."}
{"text": "Conventional means for delivering active agents are often severely limited by biological, chemical, and physical barriers. Typically, these barriers are imposed by the environment through which delivery occurs, the environment of the target for delivery, or the target itself. Biologically or chemically active agents are particularly vulnerable to such barriers.\nFor example in the delivery to animals of biologically active or chemically active pharmacological and therapeutic agents, barriers are imposed by the body. Examples of physical barriers are the skin and various organ membranes that must be traversed before reaching a target. Chemical barriers include, but are not limited to, pH variations, lipid bi-layers, and degrading enzymes.\nThese barriers are of particular significance in the design of oral delivery systems. Oral delivery of many biologically or chemically active agents would be the route of choice for administration to animals if not for biological, chemical, and physical barriers such as varying pH in the gastro-intestinal (GI) tract, powerful digestive enzymes, and active agent impermeable gastro-intestinal membranes. Among the numerous agents which are not typically amenable to oral administration are biologically or chemically active peptides, such as calcitonin and insulin; polysaccharides, and in particular mucopolysaccharides including, but not limited to, heparin; heparinoids; antibiotics; and other organic substances. These agents are rapidly rendered ineffective or are destroyed in the gastro-intestinal tract by acid hydrolysis, enzymes, or the like.\nEarlier methods for orally administering vulnerable pharmacological agents have relied on the co-administration of adjuvants (e.g., resorcinols and non-ionic surfactants such as polyoxyethylene oleyl ether and n-hexadecylpolyethylene ether) to increase artificially the permeability of the intestinal walls, as well as the co-administration of enzymatic inhibitors (e.g., pancreatic trypsin inhibitors, diisopropylfluorophosphate (DFF) and trasylol) to inhibit enzymatic degradation.\nLiposomes have also been described as drug delivery systems for insulin and heparin. See, for example, U.S. Pat. No. 4,239,754; Patel et al. (1976), FEBS Letters, Vol. 62, pg. 60; and Hashimoto et al. (1979), Endocrinology Japan, Vol. 26, pg. 337.\nHowever, broad spectrum use of such drug delivery systems is precluded because: (1) the systems require toxic amounts of adjuvants or inhibitors; (2) suitable low molecular weight cargos, i.e. active agents, are not available; (3) the systems exhibit poor stability and inadequate shelf life; (4) the systems are difficult to manufacture; (5) the systems fail to protect the active agent (cargo); (6) the systems adversely alter the active agent; or (7) the systems fail to allow or promote absorption of the active agent.\nMore recently, microspheres of artificial polymers of mixed amino acids (proteinoids) have been used to deliver pharmaceuticals. For example, U.S. Pat. No. 4,925,673 describes drug-containing proteinoid microsphere carriers as well as methods for their preparation and use. These proteinoid microspheres are useful for the delivery of a number of active agents.\nThere is still a need in the art for simple, inexpensive delivery systems which are easily prepared and which can deliver a broad range of active agents."}
{"text": "In many existing networked systems, client nodes establish connections to host nodes in a network. The client nodes and host nodes are then able to communicate for a duration of time. In many situations, an Internet Service Provider (ISP) may manage the host nodes. For example, the ISP may configure the host nodes with groups of modems contained in modem pools, and a client node may acquire access to the Internet by connecting to one of the modems in the pool. The client nodes may include home or business computing systems. In these scenarios, the host nodes may often be referred to as Points of Presence, or PoP's. In certain scenarios, the client nodes may be able to directly communicate with the host nodes using broadband connections.\nOften, host nodes in a network will maintain policy information. This policy information may include access information that may be used by the client nodes to access the host nodes. For example, the access information could include modem dial-in numbers, Internet Protocol (IP) addresses, and the like. The host nodes then provide the policy information to the various client nodes in the network. In certain scenarios, the host nodes may even provide configuration information, such as modem configuration information, to the client nodes. The client nodes may use this initial configuration information to set up their dial-up modems, and may also use the policy information to select a dial-in access number for a given host node.\nIf client nodes do not have access to local dial-in access numbers, they may often use a more generic access number to make an initial connection. For example, a client node may use an “800” access number to access a host node using a dial-up connection. The client node may then specify a particular area code, and the host node may respond by providing a list of dial-in access numbers that correspond to the specified area code. Subsequently, the client node may establish dial-up connections using one of the provided access numbers.\nIn certain situations, the client or host nodes may collect various forms of historical or statistical information relating to their network connections. For example, for modem-based connections, a client node or a host node may collect information relating to the total number of calls initiated or received, the call success rate, average session lengths, and the like. This type of information can be displayed to an administrator on a graphical user interface (GUI) to provide a summary of network activity, and can also be distributed to users in a report format."}
{"text": "The present invention relates to an adjustment device for a lateral register of at least two adjacent flexible printing plates. Axially adjustable strips that accomplish the lateral registration are spaced circumferentially.\nA lateral register adjustment device for printing plates is known from DE 42 41 588 C2. Printing plates have beveled edges which are fastened in a slit or in several slits extending in an axis-parallel direction along the circumference of the printing cylinder. Each one of the printing plates has a recess for this, which recess is engaged by a respective register pin which, in turn, is fixedly connected with one insertion strip, respectively. This insertion strip is movable in an axis-parallel direction in the slit. The insertion strips are arranged on top of each other in the slit and are each connected with a separate adjustment system.\nThe object of the present invention is based on providing an adjustment device for a lateral register.\nIn accordance with the present invention, this object is attained by the provision of printing plates whose beveled ends are received in a slit or channel in the printing cylinder. Each plate end has a register recess that will receive a register pin which is carried by an axially displaceable insertion strip. The insertion strips are situated one behind the other in the circumferential direction of the printing cylinder.\nThe advantages which can be achieved by means of the present invention rest, in particular, in that a lateral register adjustment device for printing plates with beveled edges has been created (slit clamping), whose insertion strips are secured against being ejected from the slit by simple means. Furthermore, the slit can be kept narrow, even if solid insertion strips are used. The beginning or leading end of each printing plate is provided with a lateral register adjustment device."}
{"text": "(a) Field of the Invention\nThe present invention relates to a system for detecting an amount of ink consumption or a recorded amount in an ink jet recording device, such as an ink jet printer.\n(b) Description of the Related Art\nA system for detecting a recorded amount or amount of ink consumption in an ink jet recording device is described in Patent Publication JP-A-3 (1991)-246058, for example. Briefly, the system includes a page counter for storing an amount of ink consumption which is converted from an amount of recorded data, an increment section for sequentially updating the page counter if there is an ink consumption resulting from data recording, and an alarm signal generator for generating an alarm when the residual amount of ink is considered below a threshold based on the initial amount of the ink and the accumulated ink consumption.\nSpecifically, when a new ink reservoir is placed in the ink jet printer, a signal indicating the replacement of the ink reservoir is supplied to a controller in the ink jet printer to clear the page counter. If a print command signal is received in the ink jet printer, the printer effects a recording operation by ejecting ink supplied from the ink reservoir. When an amount of data corresponding to one page is recorded by the printer, the increment section increments the count of the page counter by one page.\nIn this manner, the page counter is updated at each time the recording operation for one page is completed. If the amount of ink corresponding to the volume in the ink reservoir is counted in the page counter, the controller determines that recording operation for a specified number of pages has been completed, displays an alarm signal that the ink reservoir is empty or close to empty and stops further recording operation. A non-volatile memory is used for the page counter, thereby allowing the count representing the ink consumption or residual amount of the ink to be maintained without a power supply to the printer.\nIn the conventional ink jet printer, a user for a computer system including the ink jet printer knows the emptying of the ink reservoir only when the printing is stopped or interrupted by the alarm from the printer, which causes the computer work done by the user to be inconvenienced."}
{"text": "(1) Field of the Invention\nThe present invention relates to a process for cleaning containers, in particular containers such as casks or kegs. Particularly, the present invention relates to a process for cleaning such containers in which a cleaning solution is poured into a container, after which the container is rinsed with a suitable rinsing fluid such as water or the like. The present invention is further directed to an apparatus for cleaning containers in accordance with this process.\n(2) State of the Prior Art\nIn the beverage industry, as in many other branches of industry, large fluid containers are used. In the beverage industry, such containers are casks or kegs that can have valves mounted thereto. These types of containers are normally returnable to the supplier after use by the customer for refilling at a bottling plant. Prior to a refilling operation, however, the containers have to be cleaned in order to remove any residue from the beverage or other liquid that was previously in the container and any other contaminants. The beverage industry, being a part of the food industry overall, has particularly high demands placed thereon for the degree of cleanliness of the containers used for beverages.\nCleaning of such containers is normally performed by first pre-rinsing the container in order to remove the coarsest residue. Thereafter, the container is rinsed with a cleaning solution having a high chemical content, or a similar cleaner, so that even stubborn contaminants are removed due to the high chemical content of the solution. Of course, before each container can be refilled and reused, the container has to be completely chemical free. As such, the containers are \"post-rinsed\" with a post-rinsing fluid, usually water. This post-rinsing operation is customarily carried out over a fixed period of time that is sufficient to ensure that the container will be completely chemical free. Certain precautionary steps, however, are also taken in order to account for situations where, for example, the chemical concentration of the cleaning solution or cleaning agent may be raised due to some malfunction in the cleaning operation. These precautionary steps involve having the post-rinsing take place over a considerably longer period of time than would normally be necessary to clean the container. While this of course results in an additional consumption of water, this is accepted for safety reasons.\nNor is it possible to check the post-rinsing fluid as it flows out of the container after rinsing to determine when the fluid becomes totally free of chemicals. If such a check could be made, the post-rinsing operation could be controlled, for example by means of a pH meter. Conventional solutions that are used as cleaning fluids have concentrations of 1.5 to 2.5%, and thus a pH value of 14. However, this pH value is obtained when the cleaning solutions are rinsed out of a container to such a degree that the post-rinsing fluid draining from the container shows only a residual concentration of about 0.4%. It is not until the concentration drops below this value that the pH value will fall rapidly. When the post-rinsing operation is completed, a pH value of 7 (neutral value) is reached. However, with the known methods of measuring, it is not possible to determine the pH value at the same time as the pH is dropping in the post-rinsing fluid draining out of the container. A conventional pH measuring cell will need a considerably longer period of time in order to fall from a pH value of 14 to a pH value of 7 than is actually required for the post-rinsing operation. Thus the reaction speed of conventional pH measuring cells is inadequate to control the post-rinsing operation during the cleaning operation of the containers.\nGerman patent No. 3,424,711 uses at least one conductivity measuring probe in both an approach flow and a return flow of a cleaning agent being conveyed into and then away from an object being cleaned. A conductivity comparison element continuously processes paired readings of the probes and transmits a signal when the values of the paired readings reach a predetermined minimal difference. That is, when the conductivity measuring probe is located both in the approach flow at the entry to the working area of the cleaning facility and in a return flow at the exit from the working area, the change of the electrical conductivity of the cleaning agent as it passes through the work area can be determined. When the change falls below a pre-established minimal difference, the particular cleaning phase can be terminated and another working phase can be subsequently initiated.\nHowever, in this patent the post-rinse phase requires the addition of a certain delay or run-out time. Furthermore, the present inventor has found that using electrical conductance as a parameter for stopping the introduction of a rinsing fluid is unsuitable. Employment of electrical conductance was one of the subject matters set forth in parent application 08/382,904. However, in subsequent tests, it has been found that when using hot water as a rinsing fluid in the last rinsing phase, the electrical conductance reaches the value of the hot water after six seconds, i.e. at a first measuring step. Thus using electrical conductance in practice has been found not to provide acceptable results."}
{"text": "Metal containers are produced on metal container production lines at high speeds of, for example, 2000 containers per minute. The bodies of metal containers such as, for example, aluminum beverage cans, are formed by machines called “bodymakers”. A bodymaker machine may perform, for example, stretching of the aluminum to create the straight sides and bottom of the metal container. At a typical metal container manufacturing plant, multiple bodymakers (e.g., 20) may feed into a single production line. The containers are processed as they move down the production line.\nOne process step is to spray-coat the inside of the containers. For example, there may typically be about 8 spray guns or sprayers on a production line. Each container is coated by one spray gun. The spray-coat helps protect the inside of the containers from the material (e.g., liquid) with which the container is eventually filled.\nAnother process step is to “neck” the containers. There may typically be about 15 “necker pockets” which form the container necks. The neck of each container is formed by a single necker pocket. Necks are formed on the containers to reduce the diameter of the top portion of the container such that a lid for the container may be smaller than the overall diameter of the container.\nAnother process step is to apply a “rim coating” to the bottom of the container rim (i.e., the bottom-most part of the container upon which the container stands). Such a rim coating allows the containers to slide along the metal container production line more easily helping to prevent the containers from falling over.\nThe bottom outside portions or surfaces of metal containers are often embossed with a body maker identifier which is, for example, a numeric value identifying which bodymaker machine made the container. As a result, a defective container may be traced back to a particular bodymaker machine by looking at the bodymaker identifier.\nSimilarly, the bottom outside portions or surfaces of metal containers are often marked or painted with a color dot such that the particular color of a color dot indicates which spray gun coated the inside of the container. As a result, any problems with the internal spraying of a particular container may be traced back to the particular spray gun which sprayed the container.\nFurthermore, the rim coating on the bottom surface rim of a container is typically invisible under normal lighting conditions but is sensitive to ultraviolet light. That is, when the coated rim is illuminated with ultraviolet (UV) light, the rim will give off, for example, a blue color hue. As a result, the rim may be inspected under UV light to determine if the rim was properly coated.\nTracing back any particular metal container to a particular bodymaker machine, spray gun, rim coating machine may be done manually, after a defective container has been bumped off the production line. Such manual off-line tracing back is slow and inefficient and does not allow developing production problems to be identified quickly before too many defective containers are produced. Manually tracing back to a particular necker pocket is not typically done, if done at all.\nTherefore there remains a need in the art to more easily, efficiently, and effectively trace back defective metal containers to the source equipment (e.g., bodymaker, spray gun, rim coater, necker) which produced the defect.\nFurther limitations and disadvantages of conventional, traditional, and proposed approaches will become apparent to one of skill in the art, through comparison of such systems and methods with the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "Document encoders are used in a variety of situations to accomplish a number of different tasks, among them the encoding of checks. Checks are encoded to facilitate processing thereof against an account. The encoding is typically done at the bank or other central facility. It is often desirable, however, to encode checks with the required information at the retail level, such as at a grocery store, so as to decrease the overall processing time of those checks through the banking system. In such a situation, encoders are used as part of an overall cash settlement system and are driven by the central processor portion of the system, instead of operating as a stand-alone instrument. Currently available encoders attempting to operate in such a system context, however, have a number of deficiencies. Typically, they are not fast enough and hence do not provide sufficient document throughput. In addition, such encoders are frequently too large for convenient use at the retail level and also typically require a substantial amount of hand labor to operate. Further, such encoders have not been reliable in operation and thus are frequently maintenance intensive.\nThe present invention is a compact encoder apparatus designed to be used as part of an overall cash settlement or comparable system. The apparatus features an automatic check feeding system and a convenient check catch tray, both which are designed to be integral parts of the apparatus. A particular structural arrangement of system elements defining a check guide path results in the present invention overcoming many of the difficulties encountered with prior art encoders."}
{"text": "Device-to-device (D2D) communication has been introduced as an important enhancement to LTE-Advanced (Long Term Evolution) for cellular networks. The related standardization has been carried out in 3GPP (Third-Generation Partnership Project) to study both D2D discovery and direct communication techniques.\nFor D2D direct communications, broadcast communication in out-of-coverage scenario is prioritized for Release 12. It is the most critical task for D2D broadcast communication to support VoIP traffic and increase coverage. To increase coverage, narrow band (as low as 2 or 3 RBs) is a suitable choice for transmission of a broadcast VoIP packet. For a large channel bandwidth (50 RBs for 10 MHz bandwidth), we need to FDM (frequency domain multiplexing) broadcast transmissions from multiple UEs to fully use available spectral resources. However, it incurs half-duplex constraint and in-band emissions which eventually impact reception performance of the receiver."}
{"text": "Networking and computing hardware operates using a modular hardware structure where line cards, modules, blades, plug-ins, circuit packs, etc. (collectively “modules”) are selectively inserted into a chassis for the desired functionality. This connectivity in the chassis may be through a backplane, midplane, etc. (collectively “backplane”) with associated connectors. Conventionally, backplane connectors (“connectors”) are factory installed and not field replaceable. As network element systems increase in performance, functionality, speed, power, etc., situations arise where connectors are damaged due to wear, vibration, corrosion, bending, and the like. One particular situation for damaged connectors arises as a result of vibration induced connector damage that was discovered on the cooling fan units and associated backplane connectors. Such damage is anticipated to become more prevalent as module power and consequently fan speeds and vibrations increase. Conventionally, a connector issue in the field requires a new backplane which effectively requires a new chassis. Thus, connector replacement in-service or even in-field is currently not possible. It would be advantageous to provide in-field and/or in-service connector replacement to address the aforementioned connector issues."}
{"text": "Polycarbonates are commercially important materials possessing excellent physical and chemical properties which are useful in a wide range application from non-opaque impact resistance sheets to shaped articles. Generally, however, polycarbonates have rather low scratch resistance and are somewhat susceptible to attack by many common solvents and chemicals.\nPrevious efforts to overcome this low scratch resistance and susceptibility to attack by solvents have included lamination procedures and application onto the polycarbonate of a surface coating. Many of these prior art remedial efforts have been unsuccessful due to the incompatibility of the laminae and coating materials with the polycarbonate substrate. This incompatibility has resulted in stress cracking and crazing of the polycarbonate, crack propagation into the polycarbonate as a result of the brittleness of the coating, and a reduction of the advantageous properties of the polycarbonate such as, for example, impact resistance, tensile strength, non-opacity and elongation.\nThe prior art coatings for polycarbonates have included organopolysiloxanes, U.S. Pat. No. 3,707,397; polyester-melamines or acrylic-melamines, U.S. Pat. No. 3,843,390; and allyl resins, U.S. Pat. No. 2,332,461. These types of prior art coatings are generally applied from solutions of inert solvents and are cured to final properties by baking at elevated temperatures. The disadvantages of such systems are obvious. The heat curing requires a supply of thermal energy thereby adding to the cost of the system. Further, the thermal curing step is somewhat limited by the heat distortion temperature of the polycarbonate which is to be coated. Thus, in coating of polycarbonates, sheets of 30 mils and less generally cannot be coated and cured economically because of excessive warpage of the sheets during the thermal curing process.\nU.S. Pat. No. 3,968,305 describes a synthetic shaped article having a mar-resistant polymer surface layer integrated with the polymer surface body, said polymer surface layer consisting essentially of, in polymerized form, (a) 20 to 100 weight percent of a compound having a total of at least three acryloxy and/or methacryloxy groups linked with a straight chain aliphatic hydrocarbon residue having not more than 20 carbon atoms, (b) 0 to 80 weight percent of at least one copolymerizable mono- or diethylenically unsaturated compound. This type of a surface layer suffers from the fact that it generally has poor durability of adhesion after prolonged exposure to weathering.\nU.S. Pat. No. 3,968,309 describes a molded article of plastic having on its surface a cured film of a coating material comprising at least 30% by weight of at least one polyfunctional compound selected from the group consisting of polymethacryloxy compounds having a molecular weight of 250 to 800 and containing at least three methacryloyloxy groups in the molecule and polyacryloyloxy compounds having a molecular weight of 250 to 800 and containing at least three acryloyloxy groups in the molecule. This patent, however, also teaches that this coating must contain from 0.01 to 5% by weight of a fluorine-containing surfactant in order for the coated article to be acceptable. This patent teaches that when the coating material contains less than 0.01% by weight of the fluorine-containing surfactant, it is impossible to obtain a coated article having the requisite degree of surface hardness, surface smoothness, abrasion resistance and optical clarity. If the coating material contains more than 5% by weight of said fluorine-containing surfactant, the adhesion between a cured film of the coating material and a molded substrate of plastic is unsatisfactory.\nIt has now been found that a coating composition containing certain specific polyfunctional acrylic monomers in combination with resorcinol monobenzoate provides excellent and durable UV cured coatings for polycarbonate. Thus, the present invention provides certain acrylate ester monomer based UV-cured coatings for polycarbonates which adhere tenaciously and durably to the polycarbonate, are compatible with the polycarbonate, are mar, abrasion and solvent resistant, and maintain properties after prolonged exposure to weathering."}
{"text": "This invention relates to abrasive hones and particularly to segmented abrasive hones having a lubricating mechanism.\nA typical nuclear steam generator comprises a vertically oriented shell, a plurality of U-shaped tubes disposed in the shell so as to form a tube bundle, a tube sheet for supporting the tubes at the ends opposite the U-like curvature, and a dividing plate that cooperates with the tube sheet forming a primary fluid inlet plenum at one end of the tube bundle and a primary fluid outlet plenum at the other end of the tube bundle. The primary fluid having been heated by circulation through the nuclear reactor core enters the steam generator through the primary fluid inlet plenum. From the primary fluid inlet plenum, the primary fluid flows upwardly through first openings in the U-tubes near the tube sheet which supports the tubes, through the U-tube curvature downwardly through second openings in the U-tubes near the tube sheet, and into the primary fluid outlet plenum. At the same time, a secondary fluid known as feedwater, is circulated around the U-tubes in heat transfer relationship therewith, thereby transferring heat from the primary fluid in the tubes to the secondary fluid surrounding the tubes causing a portion of the secondary fluid to be converted to steam. Since the primary fluid contains radioactive particles and is isolated from the secondary fluid by the U-tube walls and the tube sheet, it is important that the U-tubes and the tube sheet be maintained defect-free so that no breaks will occur in the U-tubes or in the welds between the U-tubes and tube sheet, thus preventing contamination of the secondary fluid by the primary fluid.\nOccasionally, one or more of the heat transfer tubes may become internally dented due to a variety of possible reasons. One such reason for the internal denting is thought to be corrosion on the outside of the tube which presses in on the tube and causes the internal flow area of the tube to become constricted. Not only does this denting phenomenon restrict the flow through the particular heat transfer tube, but it also may lead to a failure in the tube at the point of denting which may allow the primary fluid to leak into the secondary fluid thus contaminating the secondary fluid. Of course, contamination of the secondary fluid must be prevented; therefore, such a dented tube must be repaired or replaced before the denting results in a crack in the tube. When it has been determined that a particular tube in a steam generator must be removed, the tube is cut and removed from its attachment to the tube sheet. Since corrosion products have been deposited between the tube and the tube sheet wherein the tube was positioned, it is necessary to remove the corrosion products from the hole in the tube sheet so that a new tube may be placed therein. Not only is it necessary to remove corrosion products from the hole in the tube sheet, but it sometimes is also necessary to slightly enlarge the hole in the tube sheet so that a new tube may be disposed therein and welded to the tube sheet. There are many devices known in the art for honing holes in tubular members so that the tubular member may be removed of corrosion products or may be enlarged. However, none of these hones have proved to be successful in applications to nuclear apparatus, such as in the replacement of tubes in nuclear steam generators. Because of the highly radioactive environment present in a nuclear steam generator, it is extremely important that the hone be able to remove the corrosion in a short amount of time and without failing. In addition, the tube sheet of a nuclear steam generator may have holes therein in excess of 20 inches in length that must be honed in a relatively short amount of time. Because of the high temperatures generated and because of the amount of corrosion encountered, the hones available to date were not capable of honing the holes in a tube sheet of a nuclear steam generator in a short period of time and without failing.\nIn U.S. Pat. No. 3,871,139 issued Mar. 18, 1975 in the name of Steven A. Rands, there is described a multiple-compliant-bristle, self-centering, self-sizing rotary abrasive hone made of multiple outwardly extending, flexible, extremely compliant plastic bristles having enlarged abrasive globules completely firmly attached to the outer ends of the bristles in a non-flaking and non-chipping manner. According to the patent to Rands, the hone is suitable for insertion into a recess or hole in a workpiece which is to be honed by relative movement of the abrasive hone with respect thereto. While the patent to Rands does describe a particular rotary hone, it does not teach a segmented flexible hone wherein a lubricant is conducted through the support for the bristles.\nTherefore, what is needed is a segmented flexible hone capable of self-lubrication and of honing large sections of tubular members in an environment wherein human access is limited."}
{"text": "The present invention relates in general to heat dissipation in semiconductors and, more particularly, to direct contact liquid cooling in electronic circuit module packages.\nIt is known in the art that combinations of various type semiconductors and/or hybrid circuits are required to achieve large scale electronic systems. It is desired that the packaging of these system functions be miniaturized in modular form to optimize cost, size, and weight advantages for the end user. The electronic systems often consume an appreciable amount of power and generate heat. As the level of modular system integration increases, the size of the module and the power consumption of the module correspondingly increases.\nThe electronic systems are often used at various pressure altitudes including submerged applications and high altitude applications. In some applications, conventional package cooling techniques such as convection and radiation may be inadequate to assure sufficient cooling and reliable electronic system performance. Likewise, package pressure constraints that are related to a wide range of pressure altitude variations have not been adequately addressed with present module technology to provide a universal pressure solution.\nHence, a need exists for an electronic circuit module package that is capable of reliable operation at various atmospheric and sub-surface altitudes and maintains good heat dissipation characteristics throughout its range of operation."}
{"text": "It is known that printing an acidic ink or fixer, typically a fluid containing an organic acid, together with a pigmented ink which is basic or neutral causes the pigments to destabilize or “crash” on the print media. This “crashing” effect causes improved optical density, strikethrough and edge acuity and decreased bleed.\nA problem with these acidic inks and fixers is that the acid partitions or cuts into printhead plastics and adhesives over time. This problem has been mostly ignored until recently because the expected lifetime of a printhead has been relatively short, e.g., a few years. Furthermore, printheads have been relatively inexpensive to replace.\nWays to crash pigments without an acidic ink or fixer have also been found. Rather than employing acids, multivalent metal salts (e.g. calcium nitrate) have been employed in an ink or fixer fluid. Such solutions have been found to minimize adhesive degradation, but have lead to poor pen reliability because of the problems of decap and kogation caused by the salts.\nAs printheads have become larger and more expensive and have been designed to last longer, there has arisen a need to solve the problems caused by acidic inks and fixers partitioning into printhead adhesives."}
{"text": "The netpage system provides a paper-based user interface to published information and interactive services.\nOn occasion, users of the netpage system will want to obtain copies of pages of a printed document provided by the netpage system."}
{"text": "This invention relates generally to safety devices for firearms and more particularly to a safety strap which can be applied to a handgun to maintain the hammer in a cocked position.\nAlthough single action semi-automatic handguns have a number of appealing features, the need to manually cock the hammer in order to prepare the weapon for use has detracted somewhat from the acceptance of this type of handgun. For example, persons who are unfamiliar with guns or who have difficulty in manually cocking the hammer are deterred from using the gun. In addition, the time consumed in cocking the weapon can cause a delay in firing that may have serious adverse consequences. While it is possible to cock the hammer in advance and store the gun with the hammer in a cocked position, the weapon is then fully prepared to be fired and can be accidentially discharged relatively easily. As a consequence, some persons are uncomfortable with the lack of safety resulting when the gun is stored with the hammer cocked.\nTo my knowledge, there have been no safety devices available in the past which lessen or eliminate the safety problems associated with storing a pistol in the cocked and ready to fire condition. All of the known safety devices are intended for use with other types of weapons and function to prevent pulling of the trigger or cocking of the hammer. For example, the strap arrangement shown in U.S. Pat. No. 835,349 to Deming is designed for use with a rifle and includes a rigid keeper element that acts to prevent cocking of the hammer rather than maintaining the hammer in the cocked position. All of the other devices known to me are in the form of trigger locks and similar rigid devices have a number of obvious drawbacks making them unsuitable in many circumstances.\nThe present invention has, as its principal goal, the provision of a safety device that permits single action semi-automatic handguns to be safely kept in the cocked position ready for use.\nMore specifically, it is an object of the invention provide a safety strap that can be applied to a handgun in a manner to maintain the hammer cocked to prevent the gun from firing until the strap is removed.\nAnother important object of the invention is to provide a safety strap that can be quickly and easily applied to the gun and, more importantly, quickly and easily released when the gun is needed. The strap can be removed simply by disconnecting the snap fasteners which hold it in place on the gun body.\nYet another object of the invention is to provide a safety strap of the character described that also maintains the safety of the gun on when the strap is applied to the gun. Further in this respect, the strap is constructed to accommodate a safety located on either side of the gun, and its effectiveness in preventing accidental discharge is increased accordingly.\nAn additional object of the invention is to provide a safety strap of the character described that is suited for use with guns that vary widely in size, style and construction. The strap can, with only minor variations, be adapted for use with handguns produced by virtually all manufacturers.\nA still further object of the invention is to provide a safety strap of the character described that is durable, lightweight, and simple and economical to construct.\nOther and further objects of the invention, together with the features of novelty appurtenant thereto, will appear in the course of the following description."}
{"text": "The present invention relates to cement mixtures which prevent or reduce the formation of micro-cracks and, more particularly, to an improved cement mixture of a conventional cement and a cement that is micro-encapsulated in time-release polymer coatings.\nConcrete made with conventional cements has a tendency to crack. Great care must be taken in the mixture formulation, joint design, and steel reinforcement design to minimize the chances of detrimental cracks developing. It will be appreciated that joints and steel reinforcement add substantial cost to the concrete application. Although additives such as chemical and mineral admixtures and plastic and metal fibers have often been used in concrete to minimize cracking, these can add considerably to the cost of the concrete.\nPatents of interest in this field include the following U.S. Pat. Nos. 5,639,298 (Miller et al.); 4,069,063 (Ball); 4,419,136 (Rice); and 5,728,209 (Bury et al). Briefly considering these patents, the Miller et al patent discloses a composition for controlling the setting time of cements, the composition including a sulphate component modified by the inclusion of magnesium sulphate in an amount of from about 0.5 to about 12% by weight. The Ball patent discloses a composition for controlling the setting time of cements wherein carbon dioxide is added to the cement mixture in an amount from about 0.1% to about 6.0% by weight. The carbon controls the setting time and stabilizes the hydrated products. The Rice and Bury et al patents relate to additives used with cements.\nAccording to the invention, there is provided an improved cement mixture for preventing or reducing the formation of micro-cracks in a concrete structure, the cement mixture comprising a conventional portland cement and a second cement which is micro-encapsulated in a time-release polymer coating so the second cement does not hydrate for at least a few, i.e., about three, hours after the bulk of the conventional cement is hydrated so as to provide autogenous healing of small cracks in the concrete structure.\nThe second cement, i.e., the cement which is micro-encapsulated in the time-release polymer coating, is preferably portland cement (like the first cement) or expansive hydraulic cement. The amount of cement which is micro-encapsulated in the time-release polymer coating in the mixture is preferably between 10% to 50% of the total mixture. In a more preferred embodiment, the mixture comprises 10% micro-encapsulated cement and 90% conventional portland cement.\nOther features and advantages of the invention will be set forth in, or apparent from, the following detailed description of the preferred embodiments of the invention."}
{"text": "The use of Rotary Kiln or cement kiln reactors for thermally treating plastic pellets, or chips, in particular granules, pellets or chips of crystallizable polyesters containing at least 75% of their acid units derived from terephthalic acid, orthophthalic acid, 2,6 naphthalate dicarboxylic acid or their respective diesters has been described previously in the patent application WO 2004/018541. An essential feature of WO 2004/018541 is the use of plug flow like behaviour to achieve uniformity of the properties of the granules at the exit. While WO 2004/018541 contemplates the use of baffles or internals it does nothing to teach the design of such baffles necessary to increase the mixing in the radial turning direction and yet maintain plug flow like behaviour.\nU.S. Pat. No. 3,767,601 describes a rotary kiln reactor for polyester flake with internals for good mixing. U.S. Pat. No. 3,767,601 discloses kilns with internals for both batch and continuous processes. A batch process by definition cannot have plug flow, and nothing in U.S. Pat. No. 3,767,601 indicates that the internals are specially configured for mixing while maintaining plug flow like behaviour."}
{"text": "1. Field of the Invention\nThe body relies upon a complex immune response system to distinguish self from non-self. The proper functioning of the immune system is vital for the long term health of the body.\nDeficient immune response can lead to the body's inability to protect itself from non-self matter. Excessive immune response can lead to the body's over reaction to what would otherwise be innocuous matter.\nAt times, the body's immune system must be controlled in order to either augment a deficient response or suppress an excessive response. For example, when organs such as kidney, heart, heart-lung, bone marrow, and liver are transplanted in humans, the body will sometimes reject the transplanted tissue by a process referred to as allograft rejection.\nIn treating allograft rejection, the immune system is frequently suppressed in a controlled manner through drug therapy. Immunosuppressant drugs are carefully administered to transplant recipients in order to help prevent allograft rejection of non-self tissue.\nOne drug which finds use as an immunosuppressant in the United States and other countries is cyclosporin A (CsA). CsA is a cyclic undecapeptide of the general structure: ##STR1##\nwherein all of the .alpha.-amino acid residues that form the cyclic backbone of cyclosporin A are of the L-configuration except .alpha.-amino acid 8 which is of the D-configuration. Amino acid residue 1 is derived from an unusual 9 carbon amino acid [2S,3R,4R,6E]-3-hydroxy-4-methyl-2-methylamino-6-octenoic acid. Amino acid residues 1, 3, 4, 6, 9, 10 and 11 are N-methylated on the amide nitrogen atoms of the cyclic backbone of cyclosporin A. Cyclosporin A is described in U.S. Pat. Nos. 4,117,118 (1978) and 4,396,542 (1983).\nCsA may have other useful properties such as antibiotic, anti-arthritic and anti-inflammatory activities and may find use in the treatment of other conditions such as diabetes, malaria and autoimmune diseases.\nA large number of CsA metabolites that retain the undecapeptide ring have been identified and reported (see Maurer, G.; Loosli, H. R.; Schreier, E.; Keller, B. Drug Metabolism and Disposition 1984, 12(1), 120-126, the structures, nomenclature and analytical data of the metabolites are incorporated herein by reference). It is not known what role, if any, these metabolites play in either the desired immunosuppressant activity or any unwanted adverse reactions when CsA is used for the prevention of allograft rejection.\nEven though CsA is a highly effective immunosuppressant drug, its use must be carefully managed because the effective dose range is narrow and excessive dosage can result in serious side effects. Renal dysfunction, hypertension, cardiovascular cramps, hirsutism, acne, tremor, convulsions, headache, gum hyperplasia, diarrhea, nausea, vomiting, hepatotoxicity, abdominal discomfort, paresthesia, flushing, leukopenia, lymphoma, sinusitis and gynecomastia have been observed in kidney, heart or liver transplant patients undergoing CsA treatment. Too little CsA can lead to graft rejection.\nManagement of CsA dosage involves careful control of the level of the drug present in the patient. Because the distribution and metabolism of CsA varies greatly between patients, and because of the wide range and severity of adverse reactions, accurate monitoring of drug level is considered essential.\nLaboratory methods for detection of cyclosporin have been developed. These techniques typically involve high performance liquid chromatography (HPLC), radioimmunoassay (RIA) or florescence polarization immunoassay (FPIA). See for example Wolf, B. A.; et al. Clinical Chem. 1989, 35(1), 120-124; Vernillet, L.; et al. Clinical Chem. 1989, 35(4), 608-611; Ball, P. E.; et al. Clinical Chem. 1988, 34(2), 257-260; Schran, H. F.; et al. Clinical Chem. 1987, 33(12), 2225-2229; Sanghvi, A.; et al. Clinical Chem. 1988, 34(9), 1904-1906; McBride, J. H.; et al. Clinical Chem. 1989, 35(8), 1726-1730; Quesniaux, V.; et al. Clinical Chem. 1987, 33(1), 32-37.\nEach of these techniques has certain limitations with regard to safety and complexity of the procedure and level of specificity for cyclosporins of interest. For example, HPLC requires long sample preparation and/or run times using high cost labor-intensive procedures; RIA presents the well-known hazards of handling radioactive materials; and FPIA, when based on non-specific mono- or polyclonal-antibodies, often fails to distinguish between CsA and its metabolites.\nA simple analytical method specific to selected cyclosporins is needed for use in cyclosporin treatment management.\nImmunoassay is a technique suited to the requirements of a simple method for measuring the amount of cyclosporin in a sample suspected of containing cyclosporin. However, most available antibodies capable of recognizing cyclosporins of interest also recognize and cross-react with closely related compounds such as cyclosporin metabolites. Because of this cross-reactivity, immunoassays dependent on these antibodies are less specific to cyclosporins of interest than might be desired. Cyclosporin immunoassays dependent on antibodies which cross-react with compounds other than cyclosporins of interest can be improved if the cross-reactive compounds are substantially inactivated toward recognition by the antibodies which recognize the cyclosporins of interest.\nThe methods and compositions of the present invention relate to a simple, specific immunoassay method for cyclosporins. The methods and compositions of the present invention also related to a method for inactivating cross-reactive compounds in immunoassay methods for cyclosporin.\n2. Brief Description of the Related Art\nWorld Patent Application No. 8,602,080 (1986), describes monoclonal antibodies selective for certain cyclosporins. These antibodies are prepared in response to cyclosporins bound to antigenic carriers through a modified amino acid residue such as a (D)-Lysine at position 8 or an (L)-threonine at position 2.\nEuropean patent application No. 283,801 (1988), describes a fluorescence immunoassay method and cyclosporin A derivatives used to raise antibodies. Cyclosporin A is linked to an immunogenic carrier via the amino acid side-chain at amino acid residue no. 1.\nCacalano, N. A.; Cleveland, W. L.; Erlanger, B. F. J. Immunol. Meth. 1989, 118, 257-263, describes the, presumably but not certainly, random photochemical grafting of cyclosporin A alkyl side-chains onto 4-benzoylbenzoic acid. These grafts are used to couple to immunogenic carriers for antibody preparation.\nU.S. Pat. No. 4,727,035 (1988), describes a method of immunoassay for cyclosporins. The method involves either radioiodine or fluorescence immunoassay.\nQuesniaux, V. F. J.; Tees, R.; Schreier, M. H.; Wenger, R. M.; Van Regenmortel, M. H. V. Molecular Immunology 1987, 24(11), 1159-1168, describes antibodies grown in response to cyclosporin-immunogen conjugates. The conjugates are attached to cyclosporin through a modified side chain on amino acids 2 or 8.\nU.S. Pat. Nos. 4,764,503 (1988), 4,639,434 (1987), and 4,384,996 (1983), describe cyclosporins modified at amino acid residue no. 8. These cyclosporins are hydroxylated on the amino acid side-chain and are useful as immunosuppressive drugs."}
{"text": "In general two types of local area networks (LANS) are presently available: collision based and token-ring.\nCollision based LANs are wired in \"parallel\" so that each station on the LAN is logically connected to the same cable as every other station. As a result when any one station transmits, its signal reaches all the other stations at roughly the same time, allowing for propagation delays. Each station on the LAN watches the destination address of the transmitted frame and determines if it is the intended receiver, if it is, the station copies the frame. If two or more stations transmissions overlap in time they will interfere with each other additively and the data will be lost. This is called a collision.\nStations attempt to avoid collisions by looking for an idle state on the line before transmitting. However, due to speed of light delays, a station will not always be able to immediately detect if another station is transmitting because that station's signal might not have reached it across the cable yet. The extreme case of this is when each station begins transmitting simultaneously. When any station detects a collision it will stop transmitting, pick a random time interval, wait that interval, and then attempt to retransmit.\nIn a network where a high volume of frames are being sent collisions can significantly limit the usable data bandwidth on a network.\nToken ring LANs avoid this limitation by allowing only one station on the network to transmit at any one time. Token ring LANs are wired in a logical ring with a transmit output from each station connected to a receive input on the next station. As long as no station has a frame to transmit a special frame called a token is circulated around the ring passed from station to station. When a station needs to transmit a frame it must wait until it receives the token. The station then removes the token from the ring, transmits its frame, then releases a new token. Thus, there can be only one token on the ring at a time and therefore only one station can be transmitting at a time.\nWhile logically connected in a ring, the physical wiring of token ring LAN forms a star with the transmit and receive pairs for each station, combined in a single cable, running to a central point where they are connected to a media access unit (MAU). The MAU merely provides the connection of one station's transmitter to the next station's receiver to complete the logical ring. A station that is not currently part of the ring (also referred to as being DEINSERTED from the ring) will be bypassed by relays in the MAU. When that station wishes to become part of the ring (when it INSERTS) it will provide a steady DC voltage signal to cause the relays in the MAU to switch and insert it into the ring.\nThe protocol that describes the operation of the particular token ring LANs of interest is defined by the IEEE standard 802.5. Each station on the ring implements the Media Access Control (MAC) level of the 802.5 standard in a combination of hardware and software. Acquiring the token before transmitting as described above is one of the basic operations defined by the MAC protocol. Once a station acquires the token, it may transmit data addressed to one or more of the other stations on the ring.\nEach station examines data on its receiver to determine if the data is addressed to it. If not, the data is merely buffered and regenerated through its transmitter to the next station downstream on the ring. In this way, each station acts as a repeater. If a station determines that data is addressed to itself, it will copy the data and change a \"frame copied\" bit in the data frame to indicate that the data has been copied. The data frame with the copied bit set is then sent to the next station downstream on the ring. Thus, when the station which originally sent the data eventually receives its own data frame back, it is able to determine that the data was successfully copied. At that point the originating station strips its transmitted frame from the ring and releases a new token. Now another station which has data to transmit may acquire the token and send data over the ring.\nPart of MAC protocol requires that one station assume a leadership role of \"active monitor.\" The active monitor can be any of the active stations on the ring and can change if the current station acting as active monitor deinserts itself from the ring or is otherwise unable to perform the duties of the active monitor. The active monitor performs several key functions in maintaining the MAC protocol. For instance, it watches for certain protocol violations and initiates recovery procedures if they are necessary. It also initiates a \"neighbor notification\" process by periodically sending (default of every seven seconds) an \"active monitor present\" frame which causes each station in turn to identify itself to its next downstream station. Neighbor notification allows network management software to obtain a map of the ring topology and inform each station of its nearest upstream active neighbor (NAUN) which is useful information in error recovery and trouble shooting. Most importantly, the active monitor provides the master or reference clock with which every station on the network must synchronize.\nA typical station on a token ring network is illustrated in prior art FIG. 1. A station 50 normally comprises a processor 52, monitor 54, and a network interface card (NIC) 56. It is critical that each station include a NIC 56 since the MAC code as well as the hardware which allows a station to insert into the network are normally stored on the NIC. Each NIC is connected to the MAU by a receive pair and transmit pair, in this case 60 and 64, respectively.\nAs shown in prior art FIG. 2, MAU 68 includes a plurality of relays and ports to connect stations on the token ring. For example, the transmit pair from port 72 is connected to the receive pair of port 76, and the transmit pair of port 76 is connected to the receive pair of port 80 and the transmit pair from port 80 is connected to the receive pair of port 82. Each port also contains relays including relays 84, 86, which connect to incoming transmit and receive pairs from a station's NIC. Relays 84, 86 remain connected to each other, thereby looping the transmit and receive pair from a station back on each other, until the station 50 applies a \"phantom voltage\" of +5 v to open the relays 84, 86 and insert the station 50 into the network. Until the phantom voltage is applied, the transmit and receive pairs of a port are electrically connected to each other.\nOccasionally, problems arise in the operation of a token ring. For instance, sometimes the ring will enter a \"beacon state\" in which the entire operation of the ring will be shut down until the problem is solved. One example of an error which would cause a beaconing ring is if a station on the ring has a bad NIC and is not correctly repeating frames. This could cause the token to be lost, and no station could then send data. Another problem that occasionally arises relates to the inability of a station to insert into the ring. In order to trouble-shoot problems such as these, testers have been developed to try to focus in on the cause or source of the error(s).\nA variety of tests are currently performed during a trouble-shooting procedure. Testers can currently measure phase jitter over a ring and test the continuity of cables used to connect stations in the ring. Prior art testers can also determine the speed at which the network is operating and not insert if the station is not at the correct speed. Prior art testers are also able to detect if a station, which has been unable to insert into the network, is trying to use an address that is already taken by a station on the ring. However, in order to perform these tests with prior art testers one or all of the users on the network are sometimes inconvenienced and/or the tests are not always conclusive.\nFor instance, in the case where it is desired perform tests on the network when all of the ports in the MAU are full, one of the stations must be disconnected to allow for access to the network by the tester. This situation is illustrated in prior art FIG. 3. Station 50 had to be disconnected from MAU 68 in order for prior art tester 90 to have access to the network. Depending on the tests being performed, the user of station 50 might be off the network for a substantial period of time. Also, if station 50 had been the cause of the error and it was removed from the network to make room for a prior art tester, it would not be clear to the user of the prior art tester 90 whether the problem was merely intermittent or if the removal of station 50 rectified the error.\nPrior art testers which measure phase jitter on the token ring require that the ring be rendered inoperable for the duration of the test. Depending on the size and usage of the ring, a delay of such a duration could be very undesirable if not intolerable.\nPrior art testers are also deficient in a number of other respects. In testing for a fault at the end of a conductor cable adjacent to the tester, an intermediate cable of several feet must be inserted between the tester and the cable conductor end when using time delay reflectory (TDR) to determine whether a fault exists in the cable connector end. In using TDR to identify faults along the cable conductor, prior art testers typically require a number of circuit elements to be changed in matching the tester's input impedance to the cable conductor impedance in order to avoid unwanted multiple return or reflected signals. Prior art testers also lack the ability to check and modify the speed at which they transmit frames into the token ring so that this rate conforms to the operating speed of the token ring. Prior art testers also lack certain information gathering techniques that are useful in isolating faults on the token ring. In addition, in order to perform the necessary token ring tests, prior art testers are required to be configured in a variety of ways. Prior art testers do not have the capability to perform all of these tests from one location in the network."}
{"text": "A portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent disclosure, as it appears in the U.S. Patent and Trademark Office patent files or records, but otherwise reserves all copyright rights whatsoever.\nThe present invention relates to optical focusing systems and more particularly relates to a computer-based semi-automatic focusing and assembly apparatus and method for focusing and assembling optical devices, such as laser diode assemblies, used in bar-code scanning systems.\nBar-code scanning systems are widely used in retailing, postal and parcel delivery services, manufacturing, warehousing and distribution and other industries for fast data transaction applications. The benefits of such scanning systems are countless. For example, in manufacturing, the use of bar code scanners could result in accurate raw material inventory, production routing, work-in-progress tracking, labor efficiency, inventory control, quality assurance, shipping precision, and so on.\nOne of the key elements of bar-code scanners is the laser diode assembly that provides the light source for the scanners. The focusing characteristics of the assembly have significant effects on the correct reading of bar codes. If the laser beam emitted from the laser diode assembly is not focused accurately, it will result in a weak signal level and therefore possible loss of information when it scans across the bar-code. The end result is that one has to scan several times before correct information can be captured. For bar-code scanners to be effective over a range of distances, the beam width of the laser beam has to be within tolerance at several specified distances (normally three distances corresponding to the near, mid and far fields are specified). To ensure that the beam width requirement is met and therefore high scanner accuracy is achieved, every laser diode assembly needs to be focused in high precision during production.\nCurrently such focusing adjustments are performed manually in most of the bar-code scanner production processes which limits productivity as well as focusing performance of the laser diode assemblies. The manual focusing setup typically used in manufacturing laser diode (LD) assemblies consists of the following components: the laser diode assembly consisting of a light beam emitting laser diode, a focusing lens, a lens holder with an aperture, a spring, and a diode holder; a mechanical fixture to hold the LD assembly; a beam scan head including a cylindrical rotating drum with a slit for beam width measurement, a photo detector for beam power measurement, and a level for gain adjustment; a guide rail to move the scan head for beam width measurement at different distances (the height of the scan head can also be adjusted); a beam scan block which controls the scan head and displays the beam width (the outputs of the scan block include an analog beam profile signal, a trigger signal, and a clock reference signal); a current source that supplies the current to the laser diode; and a laser power meter which is used to measure the intensity of a laser beam during the laser diode power-up stage.\nThe first step of the focusing process is to assemble the laser diode, the spring, and the lens together in the diode holder and use a clamp to hold the parts in place. Next, this assembly is placed in the mechanical focusing fixture. The clamp is lifted and the assembly is now ready for focusing. For each type of assembly, the operator is usually given three specified locations at which the beam width should be focused (i.e., near, mid and far fields). The operator\"\"s next task is to connect the current source to the laser diode and slowly increase the current to a specified value. The operator then manually slides the scan head to the nearest specified location and adjusts the orientation of the LD so that the beam is centered in the aperture of the scan head. The operator then adjusts one or more focusing knobs of the focusing fixture so that the beam width displayed on the scan head is within a specified tolerance. The process is then repeated at mid and far field locations. If the operator is successful in achieving beam widths within the specified tolerances at all locations, glue is carefully applied to the lens holder to fix the lens with respect to the LD. After the curing time of the glue has elapsed, the operator removes the assembly and starts focusing the next assembly. If the operator is not able to focus the beam to acceptable beam widths at all locations, the assembly is rejected.\nFollowing are some of the major problems associated with the manual focusing setup and procedures:\nMoving Scan Head\nOne of the major concerns of the current setup is that the operator has to repeatedly move the scan head on the rail back and forth during focusing. For mass production, the time spent on moving the beam scan head significantly reduces productivity.\nTime Consuming Focusing Check After Curing\nAfter curing, the beam widths at specified distances have to be checked again to see if they are still within tolerances. To do this in the existing manual setup, the operator has to move the scan head to different locations again, which further slows down the process.\nHuman Factor\nIn the existing setup, the focusing criteria is quite subjective with the operator making most of the decisions based on intuition and experience. In addition, the operator is not required to optimize focusing.\nWasted Time During Curing\nCurrently, the operator focuses only one LD assembly at a time. During the curing period, the operator is basically idle. This period of time can be utilized more efficiently if the operator can start focusing the next assembly. With the current setup this is not possible.\nNo Check on Beam Centering\nAnother concern is that the operator cannot accurately point the laser beam to the center of the scan head aperture due to the lack of an accurate beam centering target. Because of this pointing error, the beam width displayed on the scan block may not be an accurate indication of the true beam width at the specified clip level.\nFocusing Inconsistency\nCurrently, the operator\"\"s task is simply to bring the beam widths within the specified tolerances. No targeting of specific beam width values within the specified ranges is required. The result of this practice is that the assembled LD assemblies lack focusing consistency from assembly to assembly, which in turn may result in performance inconsistency of bar-code scanner products.\nLack of Process Tracking\nDuring the focusing process, some of the focusing parameters may change systematically due to errors in the focusing setup or systematic defects in the laser diodes and focusing lenses. Since no focusing data are recorded in the current focusing process, no process tracking can be done to detect such problems. Also, if a laser diode assembly results in a defective product when assembled into a bar-code scanner, no trouble shooting can be done because no data regarding to the focusing condition of the laser diode assembly at the time of assembling is available.\nNo Control Over the Amount of Glue Dispensed\nAfter focusing, the amount of glue applied for curing and the way how glue is applied are crucial. The operator has to be very careful in making sure that no excessive amount of glue is applied. If this operation is not carried out with care, then the focal distance may change and the focusing may go out of tolerance. In the currently process, there is no control over the amount of glue dispensed, which could cause serious problems if the operator is not careful or not skillful.\nNo Access to Beam Profile Information\nAnother concern is that the engineers do not have easy access to the beam profile information at different distances, especially that for the assembled laser diode assemblies, to perform a comprehensive analysis of the process. The beam profiles together with beam widths can provide a more complete and in-depth knowledge about the nature of the laser beam.\nManual Power-up Procedure\nThe current procedure of powering up the LD is completely manual. The operator is given a target laser power level and a corresponding range of current values. The operator is supposed to manually adjust the current supply until the specified laser power level is reached. There is no specification nor control over the speed at which the operator should follow to power-up the LD. The major concern here is that rapid power-up and excessive current supply may result in permanent damage to the LD. Another concern is that there is no continuous monitoring of the laser power in the current focusing process. Therefore, if the laser power is changed, it cannot be detected and erroneous focusing may thus result.\nAccordingly, it is an object of the present invention to overcome such shortcomings and to provide an improved method and apparatus for focusing and assembling visible laser diode devices.\nThe present invention is an apparatus and method for semi-automatic focusing and assembling of a light beam modifier, such as an optical lens, relative to a light beam source, such as a laser diode, to optimize the light source\"\"s beam characteristics at one or more desired locations along the path of the beam. The present invention utilizes one or more light beam analyzers positioned at one or more desired locations along the path of the beam to measure a beam characteristic, such as beam width, and to generate value signals based on these measurements. The value signals are sent to a host computer for comparison with predetermined optimal beam parameter ranges. The computer generates a focusing signal, based on the comparison, which is used to axially adjust the light beam modifier relative to the light source.\nPreferably, the light beam modifier and light source are held in axial alignment in a focusing fixture which is adjustable for pointing the light beam in a desired direction and for focusing the beam. The focusing station includes a stepper motor electrically connected to the host computer for receiving the focusing signal and axially adjusting the light beam modifier based on this signal. The focusing station is mounted at a fixed position on a carrier rail arrangement. The one or more light beam analyzers are slidably mounted on the carrier rail arrangement for positioning the analyzers at desired locations.\nIn a preferred embodiment, the present invention includes two focusing stations, three light beam analyzers and two beam splitters mounted on a perpendicular carrier rail arrangement. The two focusing stations are mounted perpendicularly and equidistant from a first beam splitter. The first beam splitter splits the light beam emitted from either focusing stations into two components. The first component is directed to a first light beam analyzer positioned at a desired location and the second component is directed to a second beam splitter. The second beam splitter further splits from the second light beam component a third component. The second and third light beam components are respectively measured by a second and a third light beam analyzer positioned on the rails at desired locations.\nThe spatial arrangement of the focusing stations, two beam splitters and three light beam analyzers upon carrier rails in itself allows for simultaneous measuring of beam characteristics at three locations and for alternately focusing and assembling two laser diode assemblies without the aid of a host computer. Preferably, however, the three light beam analyzers generate value signals based on their respective beam measurements and send these signals to a host computer. The host computer receives these signals and compares them to predetermined parameter ranges associated with the selected locations along the beam path. The host computer generates a focusing signal based on an optimization algorithm for the three locations and sends the focusing signal to the stepper motors for axially adjusting the light beam modifiers with respect to the light beam sources.\nTo provide easy pointing of the laser beam emitted from either focusing station, the present invention also preferably includes a video camera electrically connected to the host computer and facing one of the scan heads. The video camera provides a live image of the laser spot on a computer screen along with a reference target. Pointing of the laser beam is done either by manually adjusting the pointing mechanisms of the focusing stations or automatically through the use of a data acquisition system of the computer. For automatic pointing, the computer preferably includes a video frame grabber which digitizes the laser beam image captured by the video camera and detects its location on a scan head. Again, an algorithm which can be provided by computer software or other signal generating means, e.g., analog and/or digital circuitry (printed or otherwise), generates an adjustment signal based on the difference between the detected beam location and the known center of the scan head. The adjustment signal is sent to additional adjustment motors which activate the pointing mechanisms of the focusing fixture to center the beam on the scan head.\nThe apparatus also preferably includes a semi-automatic glue dispensing system. A valve-controlled glue dispensing system with a custom-made glue dispensing gun is used. When the operator presses a switch on the glue dispensing gun, a well-controlled amount of glue is dispensed to fix the light beam modifier relative to the light source. Once this is done, a signal is sent to the computer via the data acquisition system which initiates a curing timer to count down. The timer alerts the operator when the proper curing time has elapsed.\nIn another preferred embodiment of the present invention, the apparatus includes automatic laser power-up circuitry. The computer first powers up the light source to a pre-focusing current level for pointing purposes. After pointing, the circuitry then powers the light source to the nominal power level. The power level is monitored continuously and the laser power is displayed on the computer screen.\nIn a method for optimizing characteristics of a light beam emitted from a light source at one or more locations along the path of the light beam according to the present invention, a light beam modifier is axially aligned with a light beam source and the beam emitted from the source is continuously measured at one or more selected locations along the path of the beam. Value signals corresponding to these measurements are generated and are compared with predetermined parameter ranges. A focusing signal based on the comparison is generated and is sent to the light beam modifier for axially adjusting the modifier with respect to the light source. The focusing signal terminates when the value signals fall within the predetermined parameter ranges and the light beam modifier is fixed relative to the light source. In a preferred method, two beam splitters are positioned along the path of the light beam to split the beam into three components which, in turn, are continuously measured at three selected locations.\nAs a result of the present invention, a method and apparatus is provided for semi-automatic focusing and assembling of laser diode assemblies is provided. The benefits of which are the reduction of the time needed to focus a laser diode assembly, the reduction of moving parts in the system as well as the overall improvement of quality. With the present invention, beam profiles along with other important focusing data are automatically saved to data files. These data files provide optical engineers with an extensive database of information for detailed analysis of the focusing process.\nFor a better understanding of the present invention, reference is made to the following detailed description to be taken in conjunction with the accompanying drawings and its scope will be defined in the appended claims."}
{"text": "Foods, for example meat and meat products, are normally packed in finished consumer packages, in which instance the packaging operation of the food takes place either centrally for outward distribution of the consumer packages to supermarkets and other retail outlets, or locally on site at a point of sale.\nIrrespective of whether the food is packed centrally in the chain of distribution or locally at the point of sale, it is of critical importance that the package and its subsequent handling and storage of the packed food take place under such conditions and employing such methods that the packed food is given the very best possibility of reaching the consumer in pristine condition and freshness.\nThe rational distribution of foods on a large scale requires that the packaging operation and subsequent package handling be both simple and convenient, and obviously cost-effective as well, so as not to burden unnecessarily distribution economy. These taco requirements are contradictory inasmuch as pristine condition and freshness (so-called extended shelf life) generally require product-specific consumer packages and packaging environments, with inevitably higher costs per produced consumer package as a result, while simplicity and convenience in distribution generally presuppose standardized consumer packages from easily produced and readily available packaging materials which, while keeping distribution costs down, largely result in a deterioration in product protection properties in the produced consumer package.\nFor example, the packing of particularly oxygen gas sensitive and perishable foods, such as fresh meat takes place in that cut or sliced pieces of the meat are placed in suitable portions in correspondingly dimensioned plastic trays which, after filling at normal ambient atmospheric conditions, are closed by means of a thin plastic film which is applied over the filled plastic tray and is wrapped around the bottom edges of the tray and secured against the bottom of the tray for enveloping the tray and its contents. This conventional packing technique employing plastic trays is circumstantial in that it requires separate equipment for handling the pre-formed plastic trays and the thin plastic film, respectively. A further factor that renders this particular technique less convenient and attractive is that different types of foods generally require pre-formed plastic trays of different dimensions and geometric outer configurations. A further drawback in this context relates to the enveloping plastic film which must be adapted such that, at least for a certain guaranteed shelf life for the food in question, it ensures that the enveloped packaging atmosphere is not injuriously affected in respect of the packed food as a result of the migration or transport of gas through the plastic film. Consequently, plastic films of polyethylene and similar readily available polymer materials are not employed, since these generally lack the requisite gas barrier properties to afford the desired product protection for the packed food.\nIt is also known in the art to pack cut or sliced meat and similar foods under vacuum in flexible plastic packages. In order to give the finished package the desired combination of mechanical and physical protective properties, the material in these prior art plastic packages is constructed from a plurality of mutually bonded layers of the same or different materials (so-called multilayer materials) which, together in union, give the package the desired combination of protective properties. A packaging material of the multilayer type for vacuum packing of foods must possess particularly good gas barrier properties in order to maintain the vacuum inside the finished package during the guaranteed shelf life of the packed food, and thus requires that at least some of the material layers included in the packaging material consist of such a specifically gas-tight polymer, which is often extremely expensive. A multilayer material of the type under consideration here often requires interjacent material layers of chemical binder or adhesive in order to bond together the individual material layers with superior and lasting bonding strength, to form a well-integrated packaging material.\nAccording to yet a further known packaging technique, fresh meat is packed in a modified atmosphere in a package consisting of a multilayer barrier film. This so-called MAP technique (modified atmosphere packaging) is very similar to the above-described vacuum packaging technique, but differs from it principally in that the packaging atmosphere, instead of a vacuum, consists of a mixture of gases specifically composed to suit the product which is to be packed, such as O.sub.2, N.sub.2, CO.sub.2. The packaging atmosphere is composed such that the packed product (fresh meat) is given the desired shelf life, at the same time as the product (in particular when the packed meat is beef) must retain its consumer-attractive pristine, fresh appearance throughout its entire shelf life. In order to achieve and maintain the best possible MAP conditions within the package, this packaging technique requires the employment of at least one material layer of a barrier film (for example polyamide/polyethylene) in the multilayer structure of the packaging material. As has already been mentioned in connection with the vacuum packaging technique above, a packaging material of the multilayer type is complicated and, in addition expensive, in particular when it involves a material layer of a barrier film. The MAP packaging technique further suffers from the drawback that the packaging atmosphere, or gas mixture, which is required to achieve the desired shelf life and quality targets for the product necessitates careful product-adapted regulation and monitoring throughout the entire packaging process."}
{"text": "The present invention relates to the field of emergency signalling devices for use by individuals who may encounter emergency situations. Such situations include being lost or injured while travelling, hiking, or camping, or by being forced down in an airplane, or in similar circumstances.\nIt is known to use various types of rotary reflective devices as beacons or warning devices for attachment to vehicles such as automobiles or bicycles. Such devices are shown in various patents, i.e., U.S. Pat. No. 2,797,621 as a bicycle reflector; U.S. Pat. No. 1,798,052 as a warning signal; U.S. Pat. No. 4,120,561, as a rotary warning signal made with retro-reflective material. However, none of these prior devices are designed to transmit by reflection, an intelligible coded signal or transmit a specific emergency message."}
{"text": "Repairing damaged meningeal membranes has largely focused on implants (known as dural substitutes) which are grafted over damaged dura mater and are designed to replace and/or regenerate the damaged tissue. Current dural substitutes based on collagen matrices provide a good bioresorbable and safe substitute, compared to xenograft or allograft implants. Nevertheless, dural substitutes based on collagen matrices may display inferior mechanical properties such as short persistence, inferior tensile strength and low suture retention for use in infratentorial or spine areas. Synthetic dural substitutes may show improved mechanical and watertight properties, but generally are not absorbable, show a lack of conformability and are less easy to use.\nThe present disclosure relates to biocompatible and bioresorbable chitosan based implants having enhanced mechanical properties including increased suture anchoring strength and increased tensile strength. These enhanced mechanical properties may be achieved by controlling the pH during processing of the implant and may eliminate the need for a chemical cross-linking agent. The present matrices may be used in a variety of medical applications, including, for example surgical implants."}
{"text": "Many individuals die every year from \"cardiac electrical death\". The typical sequence is that the person has a congenital electrical problem or undelying organic heart disease for years. This can suddenly cause the lower part of the heart (\"ventricles\") to race out of control (ventricular \"tachycardia\") with an inefficient excessive rate. The ventricular tachycardia can then rapidly deteriorate into an electrical storm (ventricular \"fibrillation\") in which there is no pumping action. Unconsciousness occurs within 30 seconds and death follows in minutes.\nThere are three main therapies for such patients. The first therapy is antiarrhythmic drugs. Unfortunately, these have many side effects. The second therapy is the implantation of a defibrillator. The implantable defibrillator is designed to sense the ventricular tachycardia or ventricular fibrillation and deliver appropriate pacing pulses or a high energy defibrillation shock to the heart to restore normal (\"sinus\") rhythm. The third therapy is \"ablation\" in which the unstable heart cells, responsible for the arrhythmia, are destroyed through freezing or burning.\nBecause of the side effects of the drugs and the surgical risks involved with the implantable defibrillator or ablation therapy, patients are examined very carefully before receiving their therapies. Their propensity for ventricular tachycardia or ventricular fibrillation is determined by attempting to \"induce\" one of these conditions. A current carrying catheter is introduced through the leg and one end is moved about in the ventricle. Electrical pulses are then introduced into the ventricle in an attempt to destabilize the heart. If either a ventricular tachycardia or a ventricular fibrillation (collectively referred to as a ventricular arrhythmia) occurs, the patient is said to be \"inducible\" and the appropriate therapies are then prescribed.\nThe inducibility study is also referred to as an electorphysiological study or a PVS (programmed ventricular stimulation) study. While the inducibility study allows lifesaving therapies, it also has pain and risks of death associated with it. Hence there is a need for a noninvasive and safe means of predicting inducibility."}
{"text": "The present invention relates to a steel laminate gasket for an internal combustion engine to provide desired surface pressure around a hole to be sealed.\nAn engine for an automobile has been developed to seek high power and light weight. In order to satisfy these demands for the engine, an engine is made of aluminum alloy, to which high temperature and high pressure are applied.\nSince a steel laminate gasket is made of a plurality steel plates, the steel laminate gasket is strong against high temperature and pressure. Therefore, the steel laminate gasket has been used widely for a recent engine to which high temperature and pressure are applied.\nA cylinder head of a recent engine is generally made of aluminum alloy. In some cases, a cylinder block is also made of aluminum alloy. In case the cylinder head or cylinder block is made of aluminum alloy, a large amount of tightening pressure can not be applied to tighten the cylinder head and the cylinder block, because the cylinder head or cylinder block may deform by the tightening pressure.\nEspecially, when the gasket includes a wire ring around the cylinder hole, the cylinder head or cylinder block is liable to deform when tightened. Also, in case a large surface pressure is applied around a cylinder hole when the gasket is tightened, the cylinder head or cylinder block may deform around the cylinder hole.\nAs shown in FIGS. 1 and 2, a gasket 10 formed of an upper plate 11 with a bead 11a, a middle plate 12 and a lower plate 13 with a bead 13a is known. When the gasket 10 is tightened, the beads 11a, 13a deform to provide a strong surface pressure around a cylinder hole Hc. Therefore, unless the cylinder head and cylinder block are strong, the cylinder head and cylinder block may deform when tightened.\nIn U.S. Pat. No. 4,898,396, a gasket is formed of three plates, wherein upper and lower plates extend around a cylinder hole and have sealing mechanisms to seal therearound. In this patent, a surface pressure is formed adjacent to the cylinder hole.\nAccordingly, one object of the present invention is to provide a steel laminate gasket, in which the cylinder head and cylinder block can be tightened at desired surface pressure when the gasket is tightened.\nAnother object of the invention is to provide a steel laminate gasket as stated above, which can provide moderate surface pressure at an area slightly away from the cylinder hole to prevent deformation of the cylinder hole.\nA further object of the invention is to provide a steel laminate gasket as stated above, which can be easily and economically manufactured.\nFurther objects and advantages of the invention will be apparent from the following description of the invention."}
{"text": "It is known to create three-dimensional models of physical surfaces, such as a forest, an urban area or the terrain of a utility corridor, from a location sampling of the physical surface obtained using a three-dimensional scanning range finder. For example, such measurements may be obtained by flying over the terrain of interest in an aircraft and scanning the ground surface using a scanning laser range finder such as “LiDAR” (Light Detection and Ranging).\nThe resulting data, once processed, will take the form of a group of spaced-apart, discrete, measured surface points (a “point cloud”) representing the physical surface in a three-dimensional coordinate reference frame (typically Cartesian). The spacing between the measured surface points results from the inherent limitations in resolution of the scanning range finder, and can be significant when the physical surface being scanned is large and the details of the physical surface are fine, such as when conducting an aerial survey of an urban area that is required to capture fine architectural details of individual buildings. Moreover, commercially available three-dimensional scanning range finders for terrain measurement only measure depth and not color. As a result, the visual representation generated from an airborne three-dimensional scan of a ground surface appears as a collection of uncolored or, at best, monochrome (gray level toned) dots, which can make it difficult to resolve details that are distinguished more by color than by shape. Although photographic images can be overlaid on models based on geometric solids that are derived from the surface samples, the derivation of such models is a lengthy process that requires considerable manual intervention. While there is an emerging technique known as “flash LiDAR” that obtains both range and color information, it is currently much more expensive than conventional LiDAR and does not provide sufficient surface detail or range accuracy for many applications."}
{"text": "1. Field of the Invention\nThis invention relates generally to telecommunications, and, more particularly, to performing off-hook detection and control of power-mode.\n2. Description of the Related Art\nIn communications systems, particularly telephony such as a Plain Old Telephone System (POTS), it is common practice to transmit signals between a subscriber station and a central switching office via a two-wire, bi-directional communication channel. A line card generally connects the subscriber station to the central switching office. The functions of the line card include supplying talk battery, performing wake-up sequences of circuits to allow communications to take place, and the like. Voltage signals are processed and conditioned when being driven onto telecommunication lines.\nPOTS was designed primarily for voice communication, and thus provides an inadequate data transmission rate for many modern applications. To meet the demand for high-speed communication, designers have sought innovative and cost-effective solutions that would take advantage of the existing network infrastructure. Several technological solutions proposed in the telecommunications industry use the existing network of telephone wires. A promising one of these technologies is the Digital Subscriber Line (xDSL or DSL) technology.\nxDSL is making the existing network of telephone lines more robust and versatile. Once considered virtually unusable for broadband communications, an ordinary twisted pair equipped with DSL interfaces can transmit video, television, and very high-speed data. The fact that more than six hundred million telephone lines exist around the world is a compelling reason for these lines to be used as the primary transmission conduits for at least several more decades. Because DSL utilizes telephone wiring already installed in virtually every home and business in the world, it has been embraced by many as one of the more promising and viable options.\nThere are now at least three popular versions of DSL technology, namely Asymmetrical Digital Subscriber Line (ADSL), Very High-Speed Digital Subscriber Line (VDSL), and Symmetric Digital Subscriber Line (SDSL). Although each technology is generally directed at different types of users, they all share certain characteristics. For example, all four DSL systems utilize the existing, ubiquitous telephone wiring infrastructure, deliver greater bandwidth, and operate by employing special digital signal processing. Because the aforementioned technologies are well known in the art, they will not be described in detail herein.\nDSL and POTS technologies can co-exist in one line (e.g., also referred to as a “subscriber line”). Traditional analog voice band interfaces use the same frequency band, 0-4 Kilohertz (KHz), as telephone service, thereby preventing concurrent voice and data use. A DSL interface, on the other hand, operates at frequencies above the voice channels, from 25 KHz to 1.1 Megahertz (MHz). Thus, a single DSL line is capable of offering simultaneous channels for voice and data. It should be noted that the standards for certain derivatives of ADSL are still in definition as of this writing, and therefore are subject to change.\nDSL systems use digital signal processing (DSP) to increase throughput and signal quality through common copper telephone wire. It provides a downstream data transfer rate from the DSL Point-of-Presence (POP) to the subscriber location at speeds of up to 1.5 megabits per second (MBPS). The transfer rate of 1.5 MBPS, for instance, is fifty times faster than a conventional 28.8 kilobits per second (KBPS) transfer rate typically found in conventional POTS systems.\nTelephonic communication devices may be integrated into various electronic devices, such as computers, for example. Digital communications via data network lines may be made possible by these telephonic communication devices. Some of these devices may be remote devices where power consumption is of concern. When a telephonic communications device is integrated with these electronic devices, power consumption may become problematic. In order to monitor for communication signals, a minimum amount of power is required. Generally, the telephonic communication devices may be able to detect a communications state, such as an off hook state to determine whether full communication operation should to be implemented. However, one of the problems associated with the prior art is that the output detection is a high power operation and consumes an undesirable amount of power.\nDesigners have attempted to alleviate some of these problems by providing for a voltage change when a communication event, such as the reception of a ringing signal, occurs. However, the switches utilized in the state of the art devices to perform this function are generally only capable of changing the output voltage to accommodate for different communication events. This causes an enormous amount of power consumption to occur during the detection and operation of the telephonic communication devices. A simple voltage change does not alleviate all of the power consumption concerns. Efficiency in switching between power modes is generally less than desirable when simply performing a voltage change or switch.\nThe present invention is directed to overcoming, or at least reducing the effects of, one or more of the problems set forth above."}
{"text": "A bridge apparatus or a switch apparatus (hereinafter referred to as L2 switch), which is a relay apparatus for an Ethernet (registered trademark) frame (hereinafter simply described as “frame”), includes a database called FDB (Filtering Database) for accumulating an entry including a transmission source MAC address of a received frame, a port in which the frame is received, and other information. In the FDB, an entry is added every time a frame is received anew. This operation is called MAC address learning (or address learning).\nThe FDB is used when a received frame is relayed. Specifically, the L2 switch that receives a frame searches through the FDB using a destination MAC address of the received frame as a search key. If the destination MAC address is a learned MAC address, a port (an output port) corresponding to the MAC address used as the search key is obtained. If the received frame is transferred to the port, because another L2 switch of a transfer destination performs the same operation, the received frame finally reaches a target Ethernet (registered trademark) node (hereinafter simply described as “node”). On the other hand, concerning a received frame having a destination MAC address without an entry in the FDB, information concerning an output port is not obtained even if the L2 switch that receives the frame searches through the FDB (it is unknown to which port the frame should be transferred). Therefore, the L2 switch broadcasts the received frame to all ports other than a port in which the frame is received. This operation is called flooding. Such a frame relay operation is described in Non Patent Literature 1."}
{"text": "Metal powders possess utility in many applications. However, the ability to exploit these powders can be limited when their surfaces are coated in oxide films (and when other gases are present at the particle surface or within the body of the particles). In particular, applications making use of metal powders can be limited because of the oxygen content present in the powders and its influence upon metal sintering and other characteristics."}
{"text": "Field of the Invention\nThis invention relates to the adsorptive separation of metaxylene from a complex admixture of aromatic hydrocarbons including other xylenes. The invention more specifically relates to the recovery of metaxylene by selective liquid-phase adsorption using certain aluminosilicate zeolitic molecular sieves."}
{"text": "Recurrence of a cancer is a serious event, in particular the appearance of distant metastases, which can no longer be resected. The cure rate of patients in this situation is generally very poor. Patients with a high risk of recurrent disease are therefore given adjuvant treatment already after radical surgery, although the existence of remaining tumour cells cannot be diagnosed at this stage. Thus, there is a great need for a method to diagnose these patients, in order to be able to give them an adequate treatment from the very beginning of the post-surgical period.\nSeveral factors of prognostic significance for colorectal cancer (CRC) have been identified. Dukes' classification or staging based on the TNM classification gives good prognostic information, but still there is a great need for tests giving more detailed prognostic information especially for patients with Dukes' B and C or stage II and III cancers. Therefore the prognostic value of several molecular and genetic factors has been investigated. So far no single parameter, which allows individual monitoring of CRC patients, has been described. Cancer is often associated with a systemic chronic inflammation resulting in production of cytokines e.g. interleukin-6 (IL-6) and TNF-α and induction of acute phase reactants, such as C-reactive protein (CRP). High serum levels of CRP has been reported to correlate to poor prognosis, but still a fairly high recurrence rate is found in patients with normal serum levels of CRP, thus reducing the value of this parameter for individual monitoring of cancer patients. Therefore, scores based on several parameters, e.g. serum level of CRP and TNM-stage has been proposed. The importance of the immune system for the control and/or prognosis of CRC is suggested by the correlation between infiltration of lymphocytes and good prognosis/prolonged survival.\nCancer related immunosuppression is often associated with a systemic, chronic inflammation with increased pathological production of several cytokines, wherein cytokines are chemical mediators released by cells that affect the behaviour of other cells, (eg. IL-1βS, IL-1Ra IL-6, IL-10, IL-17, TNF-α, PGE2, TGF-β). IL-6, and TNF-α are involved in a paraneoplastic syndrome frequently found in cancer patients. This syndrome is characterized by a poor performance status, low-grade fever, anorexia, weight-loss, fatigue and distortion of various biochemical laboratory parameters. This condition is said to correlate to the tumour burden of the patient, being worse in more advanced disease. In particular, IL-6, a pleiotropic, proinflammatory cytokine, is of importance for the regulation of immune reactivity (Lotz: 1995. Cancer Treat Res 80:209; Barton B E. 1996. Med Res Rev 16:87). In addition to its regulatory role in immune and inflammatory responses, it regulates hepatic acute-phase protein synthesis, hematopoiesis and bone metabolism. IL-6 production can be induced by a large variety of stimuli which includes;                Other cytokines (such as IL-1β, IL-17 and TNF-α)        Microbial products (such as endotoxins, e.g. lipopolysaccharide from Gram-negative bacteria like Escherichia coli, formalin fixed or dried Staphylococcus aureus, mycobacterial cell wall components or synthetic constructs thereof (muramyl dipeptide) and Staphylococcus aureus anterotixins A (SEA) and B (SEB)        Mitogens, such as Phytohemagglutinin (PHA) and Concanavalin A (ConA)        Plasma/serum factors, such as IgG, IgA, immunocomplexes and fibrinogen or fibrin degradation products, e.g. D-dimer        Complement factors, such as C3a and C5a        Acute-phase reactants, such as C-reactive protein (CRP)        Extracellular matrix (ECM) components, such as fibronectin, vitronectin or proteolytic fragments or synthetic peptides thereof        Neuropeptides        \nIn addition, other pro-inflammatory cytokines (TNF-α, IL-1 and IL-8) can be induced in mononuclear cells by P-selectin (from platelets), free hemoglobin and soluble CD23. An increased serum concentration of IL-6 is often found in cancer patients, especially in patients with advanced disease and has been reported to correlate with a poor prognosis in various types of cancer, e.g. multiple myeloma, chronic lymphocytic leukemia, renal cell carcinoma, prostate cancer, ovarian cancer, metastatic breast cancer, pancreatic carcinoma and colorectal cancer. IL-6 may also be involved in the malignant process as an autocrine or paracrine growth factor in cancers like renal cell carcinoma, multiple myeloma, prostate cancer, colon cancer and Kaposi's sarcoma. Except for a possible prognostic value, the information provided by determination of serum concentrations of factors such as IL-6 can, can be questioned. The cellular origin of these serum factors is often unknown, they are immensely diluted in serum and their half-life is not accounted for. The information obtained from serum determinations of such factors is therefore most likely just a faint reflection of the mechanisms behind their appearance in serum.\nAn increased systemic concentration of IL-6 has been reported to correlate with a poor response to chemotherapy and it is well documented that patients with increased IL-6 and/or CRP serum levels cannot be successfully treated with immunotherapy.\nThe source of serum IL-6 in cancer patients is still somewhat unclear. Although IL-6 can be produced by a broad range of normal cells such as fibroblasts, monocytes, T lymphocytes, endothelial cells and keratinocytes it is generally assumed to be derived from the tumour as it may be produced by certain malignant cells. In a study of 12 colon cancer patients, IL-6 mRNA was expressed in tumour tissue in 83% of the cases but only in one patient in normal mucosa (Piancatelli et al. 1999. J Immunother 22:25) and in another study by Kinoshita et al. it was shown that the IL-6 concentration in tumour tissue of colorectal carcinoma patients were significantly higher than in normal mucosa and was correlated with the serum IL-6 concentration. In addition, serum level of IL-6 was correlated with such parameters as tumour size and proliferative activity of the tumour (Kinoshita et al. 1999 Cancer 85:2526)\nBased on these observations IL-6 was chosen as an interesting parameter of dysregulation of the immune system in cancer patients. In the present study LPS-induced IL-6 production by blood cells from CRC patients was determined preoperatively and a possible correlation to the prognosis of these patients was analysed."}
{"text": "This invention relates to arrangements for attaching a headlight or a lamp, both referred to herein as a lamp, to a vehicle part in which the lamp includes a housing with a transparent cover, referred to herein as a lens, which is attachable as a structural unit to the vehicle part by fasteners which are not covered by the lens.\nEuropean Application No. 0 604 306 A1 discloses a lamp mount arrangement in which fastening devices for the housing are covered by the lens. In this arrangement, however, the housing is first attached to a vehicle part and then the lens is mounted on the housing. The fastening devices are not accessible without removing the lens.\nGerman Patent No. 43 23 854 describes an adjusting device for a headlight which is located between a headlight and a blinker. The adjusting arrangement in this case is covered by the lens of the blinker.\nSuch covering of fasteners or adjusting devices is possible only if a second lamp abuts the lamp or the headlight. If a lamp or a headlight is fixed in a vehicle part without an adjacent lamp, that arrangement is not possible. Moreover, that arrangement requires a minimum gap width for introduction of an adjusting tool.\nIt is known in the prior art that a lamp fastening or adjusting arrangement may be covered with adjacent body parts or arranged on a rear housing wall of the lamp. Installation or actuation of such arrangements, however, is possible only if the lamp is adjacent to a removable or openable body part and accessibility is made possible by opening or removal of such a body part. If a lamp is located in the center of a body part, as for example on a fender, the fastening device cannot be covered by a hood, bumper, or the like. Furthermore, for reasons of design, it may be difficult to provide assembly-friendly access to the fastening device from the side of the fender facing the engine compartment.\nIt is also known that a lamp may be attached by one or more fastening devices, for example screws, which are accessible through the lens. In such arrangements, the fastening devices, in the absence of an additional covering, are recognizable when viewing the lamp from the front. Such arrangements are unacceptable in today's exacting aerodynamic and stylistic requirements."}
{"text": "1. Technical Field\nThis application relates to chemical analysis and more particularly to systems and methods that use nuclear quadrupole resonance.\n2. Background Information\nIt is known that an atom with more than one unpaired nuclear particle (protons or neutrons) will have a charge distribution which results in an electric quadrupole moment. Allowed nuclear energy levels are shifted unequally due to the interaction of the nuclear charge with an electric field gradient supplied by the non-uniform distribution electron density (e.g. from bonding electrons) and/or surrounding ions. This so-called Nuclear Quadrupole Resonance (NQR) effect results when transitions are induced between these nuclear levels by an externally applied radio frequency (RF) field. This electromagnetic field thus induces a magnetic resonance, unique to each material, without using a magnet. A typically NQR detection system consists of a radio frequency (RF) power source, an emitter to produce the electromagnetic excitation field, and a detector circuit which monitors for a RF NQR response coming from the object being analyzed.\nNQR has a number of practical uses, such as the detection of land mines, or of narcotics or explosives concealed in luggage, or remote monitoring of fluid levels such as in oil wells."}
{"text": "1. Field of the Invention\nThe present invention is generally directed to articles having a spinel crystal structure, and includes articles such as boules, wafers, substrates, and active devices incorporating same. In addition, the present invention relates generally to methods for forming such articles.\n2. Description of the Related Art\nActive optoelectronic devices, such as light-emitting diodes (LEDs) and laser diodes oftentimes will utilize nitride-based semiconductor layers for the active layer of the device. In this regard, the family of gallium nitride (GaN) materials, which broadly includes Ga(Al, In)N materials, have been utilized as a direct transition-type semiconductor material having a band gap that may be manipulated over a fairly wide range, on the order of about 2 to 6 eV.\nIn order to take advantage of the optoelectronic characteristics of such nitride-based semiconductor materials, they generally are formed as a single crystal. In this regard, it is generally not pragmatic to form bulk monocrystalline boules of nitride-based semiconductor material. Accordingly, the industry typically has sought to deposit such materials as a monocrystalline layer, such as by epitaxial growth, on an appropriate substrate. It is desired that the substrate on which the nitride-based semiconductor layer is deposited has a compatible crystal structure, to manifest the desired crystal structure in the as-deposited active layer. While such nitride-based materials, such as GaN and AlN can exist in several different crystal states, typically the desired crystal structure is wurtzite rather than zinc blende. In an effort to closely match the desired wurtzite crystal structure, the art has utilized monocrystalline alumina in the form of sapphire (corundum), and specifically oriented the sapphire substrate so as to provide an appropriate crystallographic surface on which the active layer is deposited. However, sapphire suffers from numerous drawbacks. For example, sapphire does not exhibit a cleavage plane that can be used to fabricate active devices. In this regard, it is generally desirable to dice the wafer into individual die (forming active devices, each having a device substrate) by cleavage rather than by slicing or sawing, as cleavage may reduce manufacturing costs and may simplify the manufacturing process.\nIn contrast, materials having the spinel crystallographic structure, if oriented properly, demonstrate a cleavage plane, the projection of which in the surface of the wafer is generally parallel to a cleavage plane of the nitride active layer, which permits predictable and reliable device fabrication. Despite the technical superiority of spinel over sapphire, a number of processing hurdles exist, resulting in somewhat limited economic feasibility. While the industry has sought to create spinel substrates by a technique known as flame fusion, the so called “Verneuil” technique, such a technique is relatively difficult to carry out, and extremely high processing temperatures have been traced to compositional inhomogeneities in the formed boule.\nThe industry has also sought to develop single crystalline spinel boules from melt-based process techniques, which include techniques such as the so-called Czochralski technique, among others. In such melt-based techniques, generally a stoichiometric crystal (typically MgO.Al2O3, having an MgO:Al2O3 ratio of 1:1) is grown from a batch melt, rather than flame-melted that involves solidification on a solid surface. While melt-based techniques have shown much promise for the creation of single-crystal spinel substrates, the process is relatively difficult to control, and suffers from undesirably low yield rates, increasing costs. In addition, extended cooling periods and annealing periods are carried out to remove residual internal mechanical strain and stress present in the boules following boule formation. Such cooling rates may be unusually low, and cooling periods significantly long, affecting throughput and increasing thermal budget and cost. In a similar manner, the extended annealing times, which may range into the hundreds of hours, further increase processing costs.\nStill further, even beyond the relatively high processing costs and despite the precautions taken in an attempt to address residual mechanical strain and stress in the crystal, oftentimes the wafers formed from boules tend to suffer from undesirably high failure rates, with frequently lower than 20% yield rates.\nIn view of the foregoing, it is generally desirable to provide improved spinel boules, wafers, substrates, and optoelectronic devices incorporating same, as well as improved methods for forming same."}
{"text": "It is known that successful in vitro culturing of differentiated cells remains especially difficult and elusive when utilizing current cell culture techniques. The lack of established and stable, normal and malignant cultures of differentiated cells has severely hampered certain fields of research, especially cancer studies and treatment.\nCurrently, the culturing of cells is done primarily with one of three methods:\nOrgan culture: maintenance of organs separated from their central vascular supply but with the organ, as an entity, left intact.\nTissue culture: culture of tissues or fragmented organs. The \"sociocellular\" relationships of the tissue architecture are preserved.\nCell culture: culture of an individual cell type divorced from other cell types.\nA marked difference in response of explanted cells occurs when they are cultured by cell rather than organ or tissue culture techniques. Methods which retain tissue architecture permit retention of tissue-specific functions including hormonal and pharmacological responses. However, the tissue normally degenerates within a few weeks, due primarily to difficulties in nutrient and waste products exchange. Cell culture procedures overcome this limitation since they use explanted tissues which are disaggregated into single cells. The cells are adpated to grow as a monolayer on a solid-state support such as, treated plastic, or as a cell suspension, and can be maintained in culture for extended periods. Nutrients are supplied by a liquid medium of defined basal composition supplemented with one of various sera. Ideally, one isolates a clonal cell line i.e., a single cell whose progeny are maintained in continuous cell culture. There is genetic uniformity, easier maintenance of the cells, and reduction of variables associated with a multicell culture system. Nevertheless, cell culture procedure usually result in distortion of cellular phenotype and karyotype; normal cells rarely adapt as permanent cell lines without developing abnormal karyotypes or losing tissue-specific functions; malignant cells adapt more easily than do benign tumor cells or normal cells; and fibroblasts or stromal components become established preferentially over epithelial cells. The difficulties of establishing differentiated cells in cell culture have been attributed to many causes. These include an inadequately defined basal medium; inadequately defined hormone requirements; the static conditions of cell culture in which nutritional and oxygen gradients develop and limit the growth and functioning of cells with strict nutritional and oxygen requirements; and loss of or damage of cell-cell junctions, perhaps essential in growth and/or differentiation or both by the cell culture procedures of mechanical and enzymic dissociation into single cell suspensions. Undoubtedly all of these have contributed to the impasse in maintaining differentiated cells in cell culture. Yet despite progress on these various fronts, the goal of routinely culturing normal differentiated cells remains elusive.\nIn a published article, namely, vol. LVIII of the Methods in Enzymology, entitled \"New Techniques for Culturing Differentiated Cells: Reconstituted Basement Membrane Rafts\" by L. M. Reid and M. Rojkind, cell culture techniques were set forth which were, in essence, attempts to simulate some of the cell-cell relationships and of the tissue matrix relevant to epithelial cells. The techniques described involved the culturing of epithelial cells on substrates of reconstituted basement membrane and in medium supplemented with hormones, serum, and with conditioned medium from feeder layers. That technique, more specifically, involved utilization of reconstituted basement membrane rafts on which were floated epithelial cells over primary cultures of mesenchymal cells normally in association with the epithelial cells. This is, therefore, not applicable to the disclosed technique which involves connective tissue-derived fibers as a substrate for body cells that are normally in contact with basal lamina material, especially differentiated epithelial cells, in cell cultures. Basal lamina material as used herein is a substance found on the surface of differentiated cells, on the basal side or surrounding said cells, and is composed of collagens, carbohydrates including glycosaminoglycans and non-collagenous proteins.\nIn a published article, namely, Vol. 2d. 17 of the Life Sciences, entitled \"A Simple, Versatile, Nondisruptive Method for the Isolation of Morphologically and Chemically Pure Basement Membranes from Several Tissues\" by Elias Meezan, J. Thomas Hjelle and Klause Brendel, a procedure was set forth for the isolation of intact basement membranes from bovine retinal and brain blood vessels, rabbit renal tubules and rat renal glomeruli. The techniques described involved a seven step procedure, with several steps utilizing primarily high concentrations of a Sodium Deoxycholate solution, concentrations much higher than in the disclosed technique. Other solutions used in the Meezan procedure contained sodium azide, sodium chloride and DNase. This procedure differs significantly from the disclosed technique is that using the Meezan procedure the resultant isolated membranes are partially denatured and not in a form usable by cells in culture. The disclosed technique, however, requires the use of functionally active, non-denatured, connective tissue-derived fibers and these, may only be acquired using the disclosed technique with its different steps, compounds and concentrations. Another significant distinction is that the disclosed method involves the use of ribonuclease in addition to DNase. The use of ribonuclease and DNase enables the production of truely pure connective tissue derived fibers, without any contaminating DNA or RNA, and this is essential for use as a culture substrate. The Meezan procedure does not use Ribonuclease. The use of DNase in the Meezan article to prevent \"viscous gel\" like DNA from interfering with the isolation procedure would not result in a pure enough or functional product required when using connective tissue-derived fibers as a culture substrate.\nIt is therefore an object of this invention to provide a method for the establishment of in vitro cultures of human or animal differentiated cells.\nIt is another object of this invention to provide a method as aforesaid in which human or animal differentiated cells are cultured using an extracellular matrix, more specifically, connective tissue-derived fibers, as a substrate.\nIt is a further object of this invention to provide a method which is applicable to both normal and malignant cells whatever their degree of differentiation.\nIt is still a further object of this invention to provide a method which significantly improves the survival and attachment rates for in vitro cultures of differentiated cells.\nIt is still another object of this invention to provide a method which significantly improves the growth properties of in vitro cultures of differentiated cells which are capable of growth in vitro.\nIt is still a further object of this invention to provide a method which will permit differentiated cells to retain a significant degree of their differentiated state.\nIt is still another object of this invention to provide a method which permits certain differentiated cells to retain much of their normal enzymatic activities.\nIt is still a further object of this invention to provide a method which permits certain differentiated cells to retain to a high degree, their ability to secrete substances such as hormones.\nIt is also an object of this invention to provide a method for the preparation and isolation of a novel culture substrate.\nIt is another object of this invention to provide tissue-specific connective tissue-derived fibers for the culturing of the related tissue-specific differentiated cells, both normal and malignant.\nIt is still a further object of this invention to provide a cell culture environment comprising a plurality of fibers as support for the cells with or without other solid state support(s) such as a petri dish or test tube, wherein differentiated cells would have an environment more nearly approximating normal tissue conditions.\nThe aforesaid objects as well as the other objects and advantages will be made more apparent in reading the following description and the adjoined claims."}
{"text": "This invention relates to a converter for producing a low voltage d.c. from a source of a.c. and more particularly but not solely to a converter for producing a low voltage d.c. from the a.c. supply voltage.\nThere are many applications where a d.c. voltage is required to be provided from an a.c. source, for example, in semiconductor circuits driven from an a.c. mains source. There are various known ways of carrying out such a conversion, for example the a.c. mains voltage may be transformed to a lower a.c. voltage by means of a transformer feeding a rectifier. This is a bulky arrangement and even if a separate winding can be employed in equipment having an existing transformer it is not always a convenient method. Another possibility is to rectify the mains voltage and reduce the high voltage d.c. so produced by means of a series voltage dropping element, e.g. a resistor or a transistor. However, such an arrangement requires a high power dissipation in the series element even if the low voltage source is required to provide a moderate supply current. Another possibility is to employ series capacitors in the a.c. supply line feeding a rectifier bridge, which is again a bulky solution.\nThe present invention has resulted from a consideration of the above mentioned disadvantages and seeks to provide a solution suitable for implementation in integrated circuit form."}
{"text": "Users are continually exposed to an ever increasing variety of clients that provide network access, such a set-top boxes, wireless phones, computers, and so on. A user of a set-top box, for instance, may view traditional television programming obtained from a broadcast network for display on a television, as well as order pay-per-view movies, receive video-on-demand (VOD), play “live” video games, and so on. Likewise, a user of a wireless phone may place and receive traditional telephone calls, as well as read email, download digital music, and so on.\nAnother such example is a digital video recorder (DVR). A DVR typically includes non-volatile storage (e.g., a hard disk) that enables the user to record desired content. DVR's also offer control functionality, such as the ability to pause content that is currently being broadcast and allows viewers to watch the content, while still in progress, from the point it was paused. The DVR plays back the content from storage, starting at the pause event, while continuing to record the currently-broadcast content. Additionally, the DVR may support other control functions, such as rewinding, fast forwarding a stored program, slow motion playback, and the like.\nTo record content using a DVR, a user was typically required to directly interact with the DVR itself. In some instances, the user could configure the DVR to record related content by specifying parameters to be matched with those of available content to locate potentially desirable content. For example, the user could specify the title of a television program so that the DVR would record each television program having that title. However, the user was not assured that the DVR would record a particular content item of interest. In other words, the user could not be certain that the potentially desirable content recorded by the DVR corresponded with the actual content the user wished to record. For example, although the DVR may be configured to record a particular television program, the DVR might fail to record a special regarding the actors in that particular television program. Therefore, when the user was located “away” from the DVR, the user could not cause the DVR to record the particular content item, even if the user had access to one or more of the clients that provide network access as previously described.\nAccordingly, there is a continuing need for improved content recordation techniques."}
{"text": "1. Field of the Invention\nThe present invention relates to computer-based systems for providing recommendations. More specifically, the present invention relates to a method and an apparatus that facilitate providing and obtaining recommendations from trusted sources.\n2. Related Art\nOne of the great advantages of the Internet is that it enables a large community of users to provide and obtain recommendations from each other. These recommendations can cover items as diverse as web sites, hotels, movies, books, and in general most types of products and services. In fact, many popular web sites, such as Eopinions or CNET, are primarily devoted to facilitating dissemination of such recommendations.\nUsers generally find recommendations provided by a large group of users to be useful in making decisions about specific items. However, recommendations for some items, such as books or movies, are highly subjective. Some users may like a particular book while other users do not. For these types of “subjective” items, it is important to be able to “trust” the source of the recommendations. For example, a user may trust a recommendation provided by his brother-in-law about a movie because the user has similar tastes in movies to his brother-in-law, whereas the user may place less trust in recommendations for the movie provided by the general public.\nUnfortunately, the extent to which a user trusts another user to make recommendations is itself highly subjective. It is consequently hard to incorporate such trust information into systems which provide recommendations."}
{"text": "Several models of relatively accessible bathtubs exist in the market, mostly with an access threshold that is lowered with respect to the ground. However, what exists in the market is suitable only for a short period of use. Indeed, such bathtubs do not adapt well with the physical degeneration of the user's age (loss of independence, need of a walker, need of a wheelchair) and become practically unusable if bathing assistance becomes necessary. In addition, people with lessening autonomy often feel embarrassed when visitors or other users of the bathroom see such a bathtub due to its appearance that is reminiscent of hospital furniture.\nOther bathtub models are better suited for people who need a walker or wheelchair or are disabled or the like. Those known bathtubs use a side door and a raised seat, but still have numerous drawbacks. For example, those known bathtubs require electricity when the bather needs to get out of the bath, so that if a power blackout occurs, these are rendered inoperative. Also, such bathtubs require a long period of time for emptying the tub. Some of the existing bathtubs have a poor design that can increase the risk of back pain for caregivers that assist the bather. Also, installing these bathtubs requires significant changes to the existing plumbing of the facility. In addition, this creates a relatively small space for future maintenance of the drain, since these are installed under the existing drain. Moreover, many of these bathtubs have replaced manual sealing with pneumatic sealing systems, to reduce the effort required by the user, such as by using compressors, which besides being very expensive, are also very noisy and run against the principle of a relaxing bath.\nAlso known are the following patent documents:\nU.S. Pat. Nos. 5,163,187 and 5,220,696 (Dannenberg et al.) assigned to Kohler Co., disclose a bathtub with a vertical axis hinged door and an inflatable seal.\nUS Patent Publication No. 2007067901 and U.S. Pat. No. 7,237,276 (Longman), disclose a bathtub with door movable along a horizontal direction between open and closed positions.\nUS Patent Publication Nos. 2011167555 and 2012198611 (Zierler et al.), assigned to Artweger GMBH & Co., disclose a closing and securing assembly for a bathtub with a water-tight door.\nWO2013125477 (Torii et al.), assigned to Sakai Medical Co., discloses lift mechanism for a bathtub with a support for a sitting section and a lift drive section which vertically moves the sitting section along the support.\nWO2015198280 (Polo Filisan), assigned to Jacuzzi Europe SPA, discloses a bathtub with a door movable between open and closed positions. The door includes two panels that are hinged together.\nU.S. Pat. Nos. 8,230,534; 9,254,066 (Torres et al.) assigned to Aquatic Co.; assigned to Kohler Co., disclose an accessible bathtub with a slide-down door that allegedly provides quick bathtub draining to minimize chilling of the bather. Indeed, a problem with known prior bathtubs is that the person inside such bathtub generally cannot leave the tub after bathing until the tub has fully drained and the door is then opened. As mentioned above, conventional tub draining systems take too long to drain, resulting in a chilling experience for the bather.\nU.S. Pat. No. 8,863,323 and Patent publication No. 2015020305 (Knapp et al.) assigned to Kohler Co, discloses a bathtub with a door movable between lowered and raised positions, door position sensors for sensing the position of the door, and a controller in electronic communication with the position sensors for controlling components of the bathtub based on the position of the door."}
{"text": "The invention relates to a brushless electric motor, in particular a miniature motor, comprising a housing, comprising at least one rotor provided with magnetized regions and mounted rotatably about a rotor axis on bearing supports of the housing, and comprising a stator having at least one stator unit, each stator unit including a set of first pole shoes, formed as claw poles, and a set of second pole shoes, formed as claw poles, which are disposed around the rotor axis, as well as a coil positioned following the rotor in the direction of the rotor axis and with its windings arranged to encircle the rotor axis, by means of which the first and second pole shoes can be magnetized.\nSuch electric motors are known in the prior art, for which the problem is to construct such motors at the lowest possible cost and as simply as possible."}
{"text": "1. Field of the Invention\nThe present invention inter alia relates to a method for producing of a single-use cell culture container.\n2. Description of the Related Art\nSuch cell culture containers are well known in the art and are at least partially made from polymer material. The known cell culture containers may be constructed from multiple layers, wherein the inner layer, i.e. the layer which gets into contact with the cell cultures, is made of a polymer material. Usually at least this layer consists of polyethylene (PE) or ethyl-vinyl acetate (EVA).\nA single-use cell culture container is a container which may be used for the biopharmaceutical or other production of peptides, proteins, antibodies or other bacterial, microbiological or cellular products in large amounts. Said process is also named fermentation and so-called bioreactors are used for this fermentation process. In the context of the described invention a single-use cell culture container, which is usually formed of a polymer material in the form of a bag for producing large amounts of biological materials, may also be called a disposable bioreactor. In the following text the term cell culture container means any shape of a single-use cell culture container or disposable cell culture container or bag or bioreactor for the production of biomaterials.\nBefore filling the known cell culture containers, which are in particular present in the form of cell culture bags, these containers are treated with gamma radiation for sterilization. This treatment is required in order to kill all microorganisms, which may disturb or affect an optimal cell growth. The sterilization is usually done such that a plurality of cell culture containers is introduced into a radiation device. Subsequently, the cell culture containers are treated with gamma radiation."}
{"text": "Piezoelectric ceramic devices produced by forming electrodes on sintered piezoelectric ceramics (hereinafter referred to as sintered ceramics) have been used in ceramic filters, speakers, ultrasonic oscillators, etc. In general, piezoelectric ceramic devices are polarized by direct current voltage so as to orient their polarization in the same direction. In order to increase the degree of polarization of the devices, high voltage is applied to the devices, and the devices are required to be resistant to the high voltage for polarization with no dielectric breakdown.\nBecause piezoelectric ceramic devices are often used under severe conditions, sintered ceramics are required to have good weather resistance, especially good moisture resistance.\nSintered ceramics are obtained by calcining a piezoelectric ceramic, then mixing it with a binder, granulating the resulting mixture, molding the resulting granular pellets and firing the molding. Usually, the molded units (\"moldings\") are stacked in a furnace to fire as many as possible. However, there is a problem in that the stacked moldings fuse so as to be combined together at high temperatures. To solve this problem, the following methods have been proposed.\nAs illustrated in FIG. 1, particles 25 are provided between stacked moldings 23. The particles used include zirconia, Al.sub.2 O.sub.3 and the like. First, the powder 25 is spread over a support (not shown). Then a molding 23 is put thereon. Further, the powder 25 is spread over the upper surface of the molding 23, and the other moldings 23 are piled up thereon one after another with spacing via the powder 25.\nAnother method is illustrated in FIG. 2 with reference to Japanese Patent Publication (JP-B) Hei-3-2821. Coarse grains 23a having the same composition as that of the moldings 23 but having a large mean grain size are dispersed in each molding 23, by which the moldings 23 are prevented from being combined together. The is referred to as a coarse grains-mixing method. The moldings 23 can be piled up without providing powders therebetween.\nIn the above-methods, piles of the moldings 23 are fired and then they are separated into individual sintered moldings 23.\nHowever, the sintered ceramics produced by the conventional methods do not sufficiently satisfy recent requests in the market for piezoelectric ceramics having much better characteristics including moisture resistance and being much low-priced.\nIn addition, the conventional method have the following problems:\n1. An additional step is required for removing powder from the sintered moldings.\n2. The sintered moldings still have traces of the powder 25.\n3. It is difficult to evenly spread the powder 25 between the adjacent moldings 23 and this causes great warping of the moldings 23. Therefore, the sintered moldings must be re-heated to remove the warping.\nThe coarse grains-mixing method is also problematic in the following points:\n1. Since the coarse grains 23a have the same composition as that of the moldings 23, some amount of the grains may be combined with the moldings 23 even under a relatively low firing temperature.\n2. Under a high firing temperature, most of the grains are combined with the moldings 23.\n3. The coarse grains 23a must be prepared in an additional step. In addition, different types of coarse grains must be prepared for different compositions of piezoelectric ceramics. Therefore, the method using such coarse grains is expensive."}
{"text": "The present disclosure relates generally to electric propulsion devices, and more particularly to such devices having improved efficiency and longer lifetimes.\nThere is an interest in efficient, high power space propulsion engines. Hall Effect Thrusters (HETs) produce thrust by ejecting ionized matter and are popular in orbit maneuvering and attitude control of many low earth orbit (LEO) and geosynchronous earth orbit (GEO) satellites.\nCurrently known HETs offer specific impulses over 2400 s, thrust over 1 N, and power exceeding 50 kW at efficiencies close to 60%. However, the commercial exploitation of Hall thrusters imposes a stringent constraint of trouble-free operation for more than 8000 hours.\nThe walls of the discharge chamber of a stationary plasma thruster (SPT) are commonly made of composite ceramic materials, for example, boron nitride, silicate oxide, and/or the like. Among many potential reasons limiting the efficiency and lifetime of a Hall thruster, an important reason is the wear of the surface layer of the discharge chamber walls. The wall erosion of the thruster occurs primarily due to plasma-wall interactions. If the ion impact energy is sufficiently large, the impact ions may cause relatively severe, undesirable sputtering of the discharge walls, the anode, and/or the hollow cathode walls. These surfaces may then develop non-uniformities (e.g. asperities) due to the sputtering, as well as to re-deposition, cracking, etc. Further, sputtered material may, in some instances, contaminate the plasma and potentially the spacecraft surface. This may significantly affect the performance of the HET, and may potentially affect the working parameter optimization.\nAlthough the lifetime issues are important to its design and potentially critical for long duration mission applications, many physical aspects in thruster plasma are yet to be understood. The lifetime of an on-board Hall thruster is expected to exceed several thousand hours. This complicates the experimental investigation and numerical prediction of the wall wear as several parameters come into play during the operational lifetime of the thruster. This generally results in a lack of reliable data on the sputtering yield under operational conditions.\nIn choosing a thruster size, one generally balances efficiency against thruster lifetime. High-energy plasma in existing technology tends to adversely interact with the walls of the thruster, as stated above. Despite significant numerical and theoretical advances of the recent past, scientists lack an adequate design to operate the Hall thruster at high power for long duration missions.\nThus, it would be desirable to provide a high efficiency and long lifetime electric propulsion device which advantageously reduces the potential for device wall erosion."}
{"text": "1. Field of the Invention\nThe invention is based on a fuel injector for injecting fuel into a combustion chamber of an internal combustion engine.\n2. Description of the Prior Art\nDE-A 103 35 340 has disclosed a fuel injector that includes a control valve for a pressure booster. The pressure booster has a working chamber that is separated from a differential pressure chamber by a pressure booster piston. The pressure change in the differential pressure chamber of the pressure booster is carried out by means of a servo valve that is activated by an associated switching valve. The differential pressure chamber contains a spring element that encompasses a lower section of the pressure booster piston, which has a diameter that is smaller than the upper section of the pressure booster piston. One end of the spring element is supported on the upper section of the pressure booster piston and the other end is supported on the injector housing. The spring element produces a restoring movement of the pressure booster piston. The disadvantage of this design is that it is not possible to pre-assemble the pressure booster piston and spring element. In addition, this embodiment has spatial disadvantages since it is necessary to select the outer diameter of the spring element to be smaller than the large piston diameter."}
{"text": "The invention relates to a welding method in which a body substantially open at opposite ends is formed from a flat metal sheet and welded along its edges. For example, roller seam welding involves forming the body with overlapping edges and feeding the region of overlap to welding rollers. The invention further relates to apparatus for carrying out the method.\nIt is known to form hollow cylindrical bodies from flat metal sheets, e.g.. can bodies or ventilation pipes, the edges being butted against one another or an overlapping region of sheet metal being formed which is then joined by seam welding. As is known, seam welding involves the bodies to be welded being fed, preceding the welding zone, through a large number of rollers adapted to match the cross-section of the body being welded, with these rollers forming a passage which surrounds the body, or a guiding tool (so-called diabolo rolls according to Swiss Patent CH-A-671904). These roller guide tools maintain the shape of the body and press its longitudinal edges into a Z-shaped rail which takes care of the butt-jointed positioning or overlapping of the longitudinal edges which is required for welding. Seam welding then takes place by advancing the metal sheet through the guiding tools with the welding rollers operating on the overlap, and normally involves the use of wire electrodes.\nIn the case of metal sheets that incorporate reliefs or recesses, or that are composed of different thicknesses of material or else have been cut out in such a way that the end surface of the rounded body is not at right-angles to the longitudinal axis, or direction of the weld seam, these conventional roller guiding tools have been known to cause problems since there is no guarantee that the workpiece will be correctly guided and centered. Therefore to date such welds have not been produced using the conventional welding machines.\nIt is therefore an object of this invention to devise a method of welding bodies, more particularly sheet metal casings, by which it is possible to weld a body formed from a sheet metal blank of any desired shape and even metal sheet incorporating recesses or differences in material thickness or hardness."}
{"text": "1. Field of the Invention\nThis invention relates generally to smoke detectors and more particularly is directed towards a new and improved, temperature-compensated, reflected light optical smoke characterized by a minimum number of components.\n2. Description of the Prior Art\nConventional optical smoke detectors typically comprise one or more light sources operating in conjunction with one or more photocells with one light source often serving as a reference for comparison with the light output of the other sources. In an obscuration type detector smoke passing between one light source and a photocell visible to the source will reduce the output of the photocell and alarm circuitry, set to a predetermined voltage level output of the photocell, will be actuated. In a reflected light type of detector, the light source is directed at a zone that is visible to the cell. If smoke is present in the zone, light reflected from the smoke will fall on the cell, increasing its output. Alarm circuitry will again be actuated in the event of such an increase.\nInsofar as many of the components in a smoke detector circuit are temperature sensitive, many detectors are provided with rather complex means to compensate for changes in ambient temperature. In general any increase in the complexity of the circuit adds to the cost of the circuit because of additional components and assembly requirements.\nAccordingly, it is an object of the present invention to provide a simple, low cost, reflected light optical smoke detector characterized by a minimum number of components.\nAnother object of this invention is to provide a temperature-compensated, reflected light optical smoke detector having selected components with complementary temperature coefficients that automatically correct for changes in ambient temperatures.\nA further object of this invention is to improve the reliability of smoke detectors by reducing the number of parts thereof."}
{"text": "This invention relates generally to communications controller and, more particularly, relates to a multiple access data communications controller for a time-division multiplex system for both voice and data communications."}
{"text": "1. Field of the Invention\nThis invention relates to electrical connectors with a ground structure for cross talk control between signal carrying conductors and, in particular, where the connector is a right angle receptacle or a right angle header.\n2. Description of Related Art\nWith the advance of technology, a high density of electronic circuits and components can be located on a printed wiring board or printed circuit board (PCB). Along with this miniaturization of electronic circuits and components, electrical connectors are needed to electrically and mechanically interconnect one PCB, such as a back panel or mother board, to one or more other PCBs, such as daughter boards. Further, it is typically desirable for such connectors to have a high signal density capacity. That is, the connectors should permit a high number of signals to pass through the connector per unit volume of the connector. However, electrical signals carried on a conductor can interfere with a signal carried on an adjacent conductor.\nThis interfering electrical effect that an electrical signal carried on a given conductor exerts on a signal carried on an adjacent conductor is referred to as \"cross talk.\" Controlling this cross talk is especially important in high density connectors. Such control can be implemented in a variety of ways.\nOne method of controlling cross talk is to connect certain terminals in a high density connector to conductive areas of a printed circuit board that are in turn grounded or connected to a predetermined ground potential. This solution is external to the connector.\nU.S. Pat. No. 4,655,518 (to Lennart B. Johnson et al.), U.S. Pat. No. 4686,607 (to Lennart B. Johnson) and U.S. Pat. No. 4,869,677 (to Lennart B. Johnson et al.) disclose a daughter board/backplane assembly with contact elements dedicated for grounding purposes. Header contact elements have contacts that can be connected to ground or a predetermined potential on a backplane. The header contact elements have other spring contacts carried by an inside header wall for touching contacts carried by a right angle receptacle outer wall. Other contacts are integral with and perpendicular to the contacts carried by the right angle receptacle outer wall for connection to the daughter board.\nU.S. Pat. No. 4,601,527 issued to Timothy A. Lemke discloses an internal shielding structure for connectors, specifically in vertical and right angle headers. The shielding structure includes a ground strip affixed to a mating surface of a header housing. The shielding structure further includes an elongated conductive spring contact with contact beams that extend in holes of side walls of the housing, lock tabs that connect to the ground strip and ground bars for connection to a grounded chassis.\nU.S. Pat. No. 4,824,383 issued to Timothy A. Lemke discloses a shielding structure in connectors or plug-type terminators for either a multiple conductor cable or a multiple tracing substrate that electrically isolates individual or groups of contact elements in the terminator to prevent or minimize cross talk between adjacent conductors and to prevent or minimize degradation of signal transmission. The terminator includes a ground structure with generally U-shaped channels. Contact elements extend into the channels. The ground structure is connected to a predetermined potential, rather than dedicating some of the contact elements for this purpose.\nU.S. Pat. No. 4,898,546 issued to Richard A. Elco et al. discloses a ground shield device for right angle connectors. A different one of the shield devices straddles alternate columns of contact elements in the connector. Each shield device clips to a tail of one of the contact elements straddled by the shield device. The shield devices are connected to ground or a predetermined potential.\nIt is an object of this invention to provide a high density right angle electrical connector for electrically and mechanically interconnecting electronic circuits and/or components controlling cross talk within the connector.\nFurthermore, it is an object of this invention to provide a high density angled or right angle electrical receptacle for electrically and mechanically interconnecting a circuit assembly and a plurality of terminals arranged in rows and columns in a header or shroud to control impedance and/or cross talk thereby to reduce or minimize degradation of signal transmission within the receptacle.\nFurthermore, it is an object of this invention to provide a high density angled or right angle electrical header for electrically and mechanically interconnecting a circuit assembly and a plurality of terminals arranged in rows and columns in a receptacle to control impedance and/or cross talk thereby to reduce, prevent or minimize degradation of signal transmission within the header."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for identifying an olfactory receptor contained in one olfactory cell.\n2. Background of the Related Art\nAn olfactory receptor is a trimeric G protein-coupled receptor (hereinafter, referred to as “GPCR”). More particularly, the olfactory receptor is one kind of trimeric G protein-coupled seven-transmembrane receptors.\nFIG. 4 shows a mechanism that a stimulus of an odor molecule to a cell membrane is converted into an electric signal.\nThe olfactory receptor is a membrane protein which is expressed on the cell membrane. The cell membrane is mainly composed of a lipid bilayer membrane. The lipid bilayer membrane has a structure of two layers each consisting of phospholipid molecules lined with a high density. This lipid bilayer membrane is shown in the center of FIG. 4 schematically. In FIG. 4, the outside of the cell is above the upper part of the lipid bilayer membrane. On the other hand, the inside of the cell is below the lower part of the lipid bilayer membrane. The trimeric G protein is placed in the vicinity of the olfactory receptor.\nThe trimeric G protein is a heterotrimer composed of an alpha subunit (Gαolf), a beta-subunit (Gβ), and a gamma subunit (Gγ). The cell contains adenylate cyclase. In FIG. 4, the adenylate cyclase is referred to as “AC”. To be more exact, the adenylate cyclase is a transmembrane-type protein. A protein RTP1S has a function to assist the olfactory receptors to be expressed in the cell membrane. It is noted that the protein RTP1S is not directly associated with the mechanism.\nNext, the mechanism is described. The odor molecule binds to the olfactory receptor. The binding separates the trimeric G protein into the alpha subunit (Gαolf) and a beta-gamma complex. The beta-gamma complex consists of the subunit Gβ and the subunit Gγ. The separated Gαolf activates the adenylate cyclase (AC). The activated adenylate cyclase (AC) converts adenosine triphosphate (ATP) into cyclic adenosine monophosphate (cAMP).\nThe cyclic adenosine monophosphate (cAMP) activates an ion channel, more particularly, for example, a cyclic nucleotide gated ion channel (CNG). The activation allows an ion to be transported from the inside of the cell to the outside of the cell, or from the outside of the cell to the inside of the cell. The degree of the transport of the ion can be measured as an electric signal.\nThe mouse contains approximately 1,300 kinds of the olfactory receptors. The olfactory receptor is referred to as “Olfr_n”. Generally, n represents a natural number. For example, an olfactory receptor having a name of “Olfr67” is known.\nOne olfactory cell includes one kind of the olfactory receptor. One olfactory cell does not include two or more kinds of the olfactory receptors.\nWhen an olfactory receptor is not expressed in a cell transfected with a gene sequence coding for an olfactory receptor, a skilled person who has acquired one olfactory cell is not able to identify the olfactory receptor included in the one olfactory cell.\nIn order to solve this problem, Non Patent Literature 1 discloses a method for identifying an olfactory receptor included in one olfactory cell among Olfr 480, Olfr 544, Olfr 545, Olfr 586, Olfr 642, Olfr 661, Olfr 672, Olfr 690, Olfr 744 and Olfr 749, using a pair of primers represented by SEQ ID: 03 and SEQ ID: 04, a pair of primers represented by SEQ ID: 03 and SEQ ID: 05, or a pair of primers represented by SEQ ID: 03 and SEQ ID: 06.\nNon Patent Literature 2 discloses a method for determining whether or not an olfactory receptor included in one olfactory cell is Olfr 16, using a pair of primers represented by SEQ ID: 07 and SEQ ID: 08.\nNon Patent Literature 3 discloses a method for identifying an olfactory receptor included in one olfactory cell among Olfr 1056, Olfr 1366 and Olfr 1484, using a pair of primers represented by SEQ ID: 03 and SEQ ID: 09."}
{"text": "1. Field of the Invention\nThe invention pertains to a coated substrate, preferably transparent, with a communication window allowing information to be transmitted therethrough.\n2. Description of the Related Art\nPatent DE 195 03 892 C1 discloses measures for reducing the shielding of coated panes in regard to microwave rays transmitting information. Such panes with electrically conducting and optically transparent coatings find their application as thermal insulation panes reflecting infrared rays (IR) and/or as electrical heating panes both for the glazing of buildings and for the glazing of vehicles.\nIn vehicles, they form together with a metal bodywork a Faraday cage, which protects the interior space of the vehicle against electromagnetic fields. In construction aloft also, premises can be protected electrically through the use of panes with an electrically conducting coating and a corresponding electrically conducting structure of the other walls. With shieldings of this nature, sensitive apparatus such as control computers can be protected, in the region of construction, against the disturbing influences due to radiophonic emitters or to radar installations of large power.\nMoreover, neither does the shielding allow through electromagnetic radiation in the microwave region, which is used as carrier wave for information. If an emitter and/or a receiver is situated alongside the antenna in a protected (vehicle) space, transmission problems occur. For example, systems for determining the position of the vehicles, for remote control, for identification, for fixing tolls, or the like are disturbed.\nIt is known for systems of layers to be structured a posteriori, mechanically or thermally, by removing lines of the layer deposited initially in a continuous manner. In particular, extraordinarily narrow slots can be produced in the layer with laser rays. In the aforesaid state of the art, there is devised as remedy in the electrically conducting layer at least one slot forming a radiating slot, with a length tuned to the wavelength of the microwave radiation and a very small free surface area, through which the energy radiating in the microwave range picked up by the conducting layer can again be decoupled. A so called communication window is thus formed.\nIf the working frequency for the transmission of information equals for example 5.8 GHz, as is envisaged for the automatic fixing of tolls on motorways (“remote payment”) (DSRC standard with a mean frequency of 5.8 GHz with circular polarization) and if the slots are provided mainly for the transmission of the microwaves of this frequency, they will be appropriately designed for the resonant length of λ/2, taking account of the dielectric constants of the glass.\nShould the information be transmitted with circularly polarized microwaves (that is to say the instantaneous plane of oscillation of the waves rotates about its axis of propagation, so that the waves oscillate inside a circular envelope curve), appropriate provision is made for gaps in the form of cruciform slots in the layer. The length of the two slots is again tuned in an appropriate manner to the wavelength of the microwaves used and corresponds to the value μ/2 of the microwaves used, taking account in a corresponding manner of the dielectric constants of the glass.\nComparative measurements of the damping of a microwave radiation of frequency 5.8 GHz indicate, in this state of the art, that with a laminated pane exhibiting radiating slots in the coating, appreciably weaker damping of transmission, for high frequency radiation, is obtained than with a fully coated laminated pane, and that it is possible to approach the damping of an uncoated laminated pane.\nDocument DE 198 17 712 C1 describes a similar substrate with a communication window in a coating, which is produced by the creation of a structure in the form of fine lines or fine patterns in a limited part of the surface of the coating and which is visually very discreet.\nThe formation of resonating structures in the conducting coating poses problems however. Results of trials have shown that the high frequency currents necessary for compensation, in the overall system formed by the dielectric substrate and the conducting coating, were not able to flow on account of the high surface resistance or of loss of the usual conducting coatings.\nDocument WO 00/72635 A1 describes a transparent substrate with a coating that reflects IR and a communication window made by the surface removal or the omission of the coating. In contradistinction to communication windows made solely in the form of fine lines, discussed in the introduction, this variant forms a disturbance of the coating which is markedly perceptible to the eye through a difference in color at the boundary of the coating.\nThis disturbance poses a problem in particular, when the coating is used at the same time for the electrical heating of the substrate. For this purpose, a voltage is applied to the coating with the aid of at least one pair of electrodes (in the form of bands), the currents having to be introduced and distributed as uniformly as possible in the surface of the layer. For vehicle panes, which are appreciably wider than high, most of the time the band shaped electrodes are situated along the longer sides of the pane, so that the heating current can flow over the shortest path over the height of the pane. At the same time, the communication windows are situated on the upper rim of the pane and here extend over a width of several centimeters.\nManifestly, each communication window affecting the homogeneity of the coating forms a disturbance of the current flow. Local temperature spikes (“hot spots”) appear which may lead to damage to the substrate (thermal stresses) and to the coating itself. This is not only the case when the coating is of large extent, but also when the communication window is formed by a larger or smaller number of discontinuous individual slots. The latter also form in the zone of the surface considered an appreciable increase in the resistance of the layer and also in addition cause the appearance of the hot spots mentioned above.\nThe document cited last proposes, as a measure aimed at reducing the disturbing effect of the large communication window, that provision be made on the edge of the latter for an electrically conducting band, which exhibits a much smaller ohmic resistance per unit area than the heating layer. This ought to cause the currents to bypass the cut-out. Preferably, a communication window is entirely flanked by such a band. The band may be fabricated by printing and baking a conducting silk screen printing paste, containing silver. It may, however, also be applied by depositing an electrically conducting lacquer or by laying a metallic band. In all cases, an electrically conducting functional bond is naturally necessary between the band and the coating.\nThe band may be masked from view by the superposition of an electrically nonconducting, opaque masking band, made for example of black enamel. Such masking bands are as a general rule composed of a nonconducting substance, black in color, which may be baked (silk screen printing paste). An infrared radiation is not reflected, but absorbed, by this substance."}
{"text": "The bacterial signaling molecules cyclic diguanosine monophosphate (c-di-GMP, 6a) and cyclic diadenosine monophosphate (c-di-AMP, 6b) are both key second messengers (Gomelsky, M., cAMP, c-di-GMP, c-di-AMP and now cGMP: Bacteria use them all, Mol. Microbiol. 2011, 79, 562-565; Römling, U., Great times for small molecules: c-di-AMP, a second messenger candidate in bacteria and archaea. Sci. Signal. 2008, 1, pe39). Because of their critical functions in essential bacterial pathways, c-di-GMP and c-di-AMP have major implications for human health. However, in spite of advances in the identification of proteins and RNA to which c-di-GMP and c-di-AMP bind, a great many additional receptors remain unknown (Galperin, M. Y.; Higdon, R.; Kolker, E., Interplay of heritage and habitat in the distribution of bacterial signal transduction systems. Mol. BioSyst. 2010, 6, 721-728). New tools are needed to determine the roles c-di-GMP and c-di-AMP have in human health."}
{"text": "The present invention relates to a new and distinct cultivar of Brachycome plant, botanically known as Brachycome segmentosa Hybrid and referred to by the cultivar name `Jumbo Misty Pink`.\nThe new Brachycome was discovered by the Inventor in a group of seedling progeny from a cross of unidentified selections of Brachycome segmentosa Hybrid in an outdoor area in Wandin North, Victoria, Australia.\nAsexual reproduction of the new Brachycome by terminal cuttings taken in a controlled environment in Wandin North, Victoria, Australia, has shown that the unique features of this new Brachycome are stable and are reproduced true to type in successive generations."}
{"text": "This invention pertains to the design of a high throughput electron beam lithography system based on an electron beam column equipped with a novel photocathode, which photocathode has a pattern composed of a periodic array of sub-wavelength apertures with a specific geometry for creating a transmission resonance.\nAs encapsulated by Moore\"\"s Law, there is a longstanding trend in the semiconductor industry toward higher device densities and correspondingly smaller device geometries. One technique for coping with these decreasing device sizes is electron beam lithography in which electrons generated by an electron source are accelerated by an electric field and focused by electron optics onto a semiconductor substrate in order to expose an electron-sensitive resist coating on the substrate surface. Electron beam lithography has the advantage of being able to achieve much higher spatial resolution than light-based lithography.\nIn a conventional electron beam lithography system, a single beam in a single column is used to expose the resist and create the desired pattern. The throughput of a conventional electron beam lithography system is limited by the total beam current and data delivery rate. As the total beam current is increased, repulsive electron-electron interactions in the beam cause excessive blur, resulting in a degradation of the resolution of the written pattern. Increasing the length of the electron column also tends to increase the blur caused by the electron-electron interactions. Several approaches are known to reduce electron-electron interactions and the associated beam blur. When the total current is distributed evenly among several columns, a reduction in the beam blur of each column results in higher resolution, and if all the columns are operated in parallel, the throughput is not compromised. This is impractical for a conventional column due to its overall dimensions and large footprint. (in addition, the data delivery rate per column is reduced. For a given throughput, multiple beam approaches reduce the incrementing rate by the number of beamlets used.) In a single electron beam column, the total current can also be divided into several beamlets, which approach also reduces the blur due to electron-electron interactions and therefore allows a higher total beam current.\nElectron lithography systems may also be divided into those in which one or more electron beams are selectively both modulated and moved over a semiconductor substrate in order to xe2x80x9cdirectly writexe2x80x9d a circuitry pattern in a serial manner without a mask, and those in which an electron emission pattern is created and electron-optically imaged onto a relatively large area of the substrate in order to expose a circuitry pattern in a concurrent manner. Such electron emission patterns may be created by illuminating a photocathode through an optical mask or by using a photocathode that incorporates a mask pattern into its own photoemmissive structure. In the alternative, the electron emission pattern may also be defined by transmitting the electron beam through a mask or reticle.\nWhen photocathodes are used to generate one or more electron beams for use in electron photolithography, it is desirable to be able to fabricate very small emission sites, since this reduces the required demagnification ratio. The optimum demagnification ratio is equal to the ratio of the source angle and the image angle at the substrate. Typical values are 2.5 mrad for the source angle, and 10 mrad for the image angle, which implies an optimum demagnification ratio of xc2xc. Given a spot size of 50 nm at the substrate, and optimum demagnification ratio of xc2xc, the size of the emission site should preferably be approximately 100 nm. Here we assume that the demagnified image and aberrations sum to the total spot size. For larger demagnification ratios, a part of the beamlet current must be cut out by the beam-limiting aperture, and therefore only a fraction of the emitted photoelectrons can be utilized for beam exposure. This requires a larger photoemission current at the photocathode, which increases blur induced by electron-electron interactions and therefore limits achievable throughput. Preferably, the size of the photoemission sites should be less than or equal to 100 nm. However, transmission through apertures substantially smaller than the shortest wavelength of practically available lasers (xcx9c250 nm) is very low. This means that the available photoemission current becomes negligible for such small aperture sizes.\nThe present invention involves a novel technique using a multiple electron beam column with a wavelength-period matched photocathode pattern for electron beam lithography. The transmission of this periodically patterned photocathode is anomalous, since it acts as an active element with a normalized transmission efficiency larger than 1; i.e., the light transmitted through the aperture is greater than the light directly impinging on it. (As used herein, xe2x80x9clightxe2x80x9d is not intended to be limited to wavelengths in the visible range.) This allows for the efficient illumination of sub-100 nm emission sites, therefore increasing the current available in each individual beamlet. As such, this arrangement allows high throughput to be achieved with very high spatial resolution lithography by simplifying and shortening the column, which reduces significantly the beam blur due to electron-electron interactions.\nIn a preferred embodiment, the invention is a method and apparatus for implementing sub-100 nm electron beam lithography at high throughput, and with an electron beam lithography column having a simplified design resulting in shorter column length, reduced electron-electron interactions and higher beam current.\nIn a class of embodiments, the invention is an electron beam lithography system based on an electron beam column having a novel photocathode, which photocathode has a pattern composed of a periodic array of apertures with a specific geometry. The spacing of the apertures is chosen so as to maximize the transmission of a laser beam through apertures significantly smaller than the photon wavelength. The patterned photocathode is illuminated by an array of laser beams to allow blanking and gray-beam modulation of the individual beams at the source level by the switching of the individual laser beams in the array. Potential applications for this invention include direct writing of a circuitry pattern on a target (e.g., a semiconductor wafer) using an electron beam, and mask patterning. In some applications, the electron beam is preferably a high voltage (50-100 kV) electron beam column; however the invention is also applicable to microcolumns operating at lower beam energies such as 1-5 kV."}
{"text": "The present invention relates to a new and improved construction of mixer arrangement and, in particular, relates to a vat or bucket mixer equipped with a mechanical drive arrangement, especially for the commingling or mixing together of liquid- and powder components, for instance when producing printing inks or the like, and which mixer is of the type comprising a rotatable mixing tool extending to the region of the inner wall of the vat and a toothed disk driven by a shaft, this toothed disk constituting a so-called dissolver disk.\nIn order to mix ttogether liquid- and powder components there are placed upon the mixer, among other things, the following specific requirements:\n(A) Disintegration or tearing apart of the agglomerates;\n(B) Optimum dispersion of the constituents; and\n(C) Good wetting of the individual pigments.\nThe properties of the different components should not be destroyed, something which oftentimes is a problem in the case of products which are difficult to work with, such as for instance heat-sensitive products, viscous pastes or thixotropic products.\nMoreover, the working operations should be accomplished in an environmentally safe manner, i.e. there should be developed as little dust as possible.\nThe heretofore known previously described mixers possess two separate drive units. Consequently, it is possible to initially carry out the pre-mixing or starting mixing operation at a relatively low rotational speed with the aid of the mixing tool, so that the power requirements are maintained small and also the formation of dust is low.\nAfter the pre-mixing there is connected the dissolver disk by means of the separate motor in order to disintegrate the agglomerates and to achieve optimum dispersion and use of the powder components. The pre-mixing phase or starting mixing phase, lasts only about five minutes, during which time there is only in operation the mixing tool, and the second phase, during which time there are in operation the mixing tool and the dissolver disk depending upon the strived for quality of the product, lasts up to one hour and more.\nBy separating the drive for the mixing tool and the dissolver disk it is possible to easily obtain variations in the rotational speed. The determination of the power requirements must be carried out individually both for the mixing tool as well as for the dissolver disk, yet in both instances is difficult. Since it must be presupposed that both demand increased power requirements, it is necessary to design both motors so as to also have the corresponding safety factor, resulting in relatively large structural requirements, and only during part of the fields of application are there in fact realized the advantages."}
{"text": "The disclosed subject matter relates generally to subsea hydrocarbon production and, more particularly, to an optical leak detector for detecting material leakage from subsea equipment.\nTo control a subsea well, a connection is established between the well and a monitoring and control station. The monitoring and control station may be located on a platform or floating vessel near the subsea installation, or alternatively in a more remote land station. The connection between the control station and the subsea installation is usually established by installing an umbilical between the two points. The umbilical may include hydraulic lines for supplying hydraulic fluid to various hydraulic actuators located on or near the well. The umbilical may also include electrical and or fiber optic lines for supplying electric power and also for communicating control signals and/or well data between the control station and the various monitoring and control devices located on or near the well.\nHydrocarbon production from the subsea well is controlled by a number of valves that are assembled into a unitary structure generally referred to as a Christmas tree. Christmas tree and wellhead systems have the principle functions of providing an interface to the in-well environment, allowing flow regulation and measurement, and permitting intervention on the well or downhole systems during the operational life of the well. The actuation of the valves in the Christmas tree is normally provided using hydraulic fluid to power hydraulic actuators that operate the valves. Hydraulic fluid is normally supplied through an umbilical running from a remote station located on a vessel or platform at the surface. Alternative systems using electrically based actuators are also possible.\nThe subsea equipment includes many possible leakage paths, such as valves, pipe junction actuators, flanges, pipe connectors, jumpers, seals, etc. The detection of leaks of hydrocarbons, hydraulic fluids, tracers and other chemicals from underwater structures is an important requirement in enhancing the environmental and operational efficiency of underwater systems such as Christmas trees and subsea processing systems. A range of underwater leak detection systems have been developed including those based on acoustic, fluorescence, temperature and gas measurements. Often these are designed to be used as part of survey operations, but in a few instances as permanent monitors. Such devices are typically portable and are mounted to mobile devices, such as remotely operated vehicles (ROVs).\nAcoustic devices are capable of detecting leaks from a wide area via the noise that may be produced by material leaking from underwater structures. Such acoustic systems detect only the secondary effect of the leak (i.e., the noise), hence, the application of acoustic detectors is significantly restricted in noisy environments. Further, acoustic detectors are not generally able to accurately locate a leak.\nThis section of this document is intended to introduce various aspects of art that may be related to various aspects of the disclosed subject matter described and/or claimed below. This section provides background information to facilitate a better understanding of the various aspects of the disclosed subject matter. It should be understood that the statements in this section of this document are to be read in this light, and not as admissions of prior art. The disclosed subject matter is directed to overcoming, or at least reducing the effects of, one or more of the problems set forth above."}
{"text": "The invention relates to a spray head for a spraying tool, especially for spraying mold halves in a pressure casting machine. The spraying tool has spray nozzles fed with spraying agents and is fixed upon spray plates or spray bars which are assembled in a modular manner.\nA spray head of this kind is disclosed in U.S. Pat. No. 4,760,961. The elements therein, equipped with spray nozzles, are assembled together column-wise in vertical or horizontal lines, and the media to be sprayed can be passed from element to element. These individual spray elements and also spacing parts inserted between them are held together by threaded bolts which are passed through openings in the elements. This method has the disadvantage that, for different sizes of the assembled spray elements or plates, different fasteners of various length have to be prepared. It is also disadvantageous that, in the case of a plurality of module-like elements arranged in a row, the threaded bolts or tie bolts have to be released and with them all the extensions that they hold together. Under certain circumstances, if an additional extension or reduction of elements not threaded on at the end is needed, but at the beginning, i.e., close to the distribution bar, all of the elements have to be taken off from the tie bolt. In this case dirt can enter at the points of separation.\nEP 0 724 486 B1 discloses a spray bar on which plate-like elements are assembled. There too, however, the component parts are held together by threaded bolts or tie bolts of appropriate length.\nAn object of the invention is to design a spray bar such that a simple modular assembling of individual spray elements or connecting parts is possible without the need for fastening the parts together with threaded bolts or the like.\nTo achieve this and other objects, a spray head has coupling elements. The coupling elements have passages for the media to be sprayed, such as air, control means and spray media, are provided on the ends of the parts to be joined together, and are connected to the channels supplying the media. The coupling elements reach into the two chambers of the ends to be joined and are provided with hollow spaces with at least one sloping plane, which are engaged by clamping means which can be operated from the outside and force the parts tightly against one another.\nBy this configuration each element is provided on its connecting surfaces with connecting means which permit the elements to be released at any point, without the need to thread bolts or the like through holes.\nIn an advantageous configuration of the invention the coupling elements are provided at both face ends with sealing rings, so that they also provide the necessary seal between two elements which are to be assembled together.\nThe coupling elements can be made simply as cylindrical pins having two recesses each. The recesses are provided with planes sloping up toward one face surface. In this case the clamping means can be configured simply as screws with a conical tip which cooperate with the recesses.\nIn further development of the invention, pairs of chambers can be provided parallel to one another at each end of the component elements, and these chambers can be coaxial with the passages carrying control air and parting agents. Also at least one passage can be provided for drying air or compressed air between the passages provided for control air and parting agents.\nThe component elements can be provided with angled end portions for the connection of additional component elements or plates at an angle, and the end portions of attached corner pieces are formed in a simple configuration with angles of any kind, e.g., 45° or 90°. With this configuration there is also the possibility that component elements can be provided, for example, in bar or plate form turned away from one another at different directions and angles. Thus any desired spatial arrangement of spray elements is possible. To lengthen or shorten the spray bars or spray plates it is necessary only to separate them at the point at which the lengthening or shortening is to be made. All of the coupling elements in the entire spray head are in that case identical.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "1. Technical Field\nAspects of the present invention relate to an image forming device.\n2. Related Art\nMulti-color image forming devices (e.g., a color laser printer) having a plurality of chargers, the number of which is equal to the number of colors of developer, have been used. Among the image forming devices of this type, an image forming device, in which a high-voltage unit (a high voltage circuit) for applying a high voltage is shared by the plurality of chargers so as to reduce the number of parts and to make the image forming device compact in size, has been used."}
{"text": "The methanol production process generally involves directing a compressed synthesis gas comprising hydrogen, carbon monoxide and carbon dioxide at an elevated temperature and pressure to a methanol converter reactor containing one or more beds of a methanol synthesis catalyst such as a copper and zinc oxide catalyst. The carbon monoxide and carbon dioxide in the synthesis gas react with the hydrogen to form methanol across the catalyst. The methanol synthesis process is usually operated in a loop where a portion of the compressed synthesis gas is converted to methanol each pass through the methanol converter reactor. Most of the unconverted gas is recycled to the methanol converter. A small portion is purged to prevent the buildup of inerts such as nitrogen, argon and methane. Methanol product is recovered by cooling the methanol product gas stream to a temperature below the dew point of the methanol such that a product composition comprising crude methanol and water condenses out, with the remaining gas being recycled through the methanol converter reactor. The crude methanol and water produced in the methanol converter reactor are typically reduced in pressure in a let-down or “flash” vessel. Since most crude methanol contains a range of impurities, including higher alcohols, the crude methanol must be purified so as to remove such impurities to produce methanol of chemical grade quality. The preferred technique used for methanol purification is a distillation process.\nSynthesis gas used for methanol synthesis is typically characterized by the stoichiometric ratio (H2−CO2)/(CO+CO2), often referred to as the module or stoichiometric number, wherein H2, CO2 and CO denote the mole fractions of hydrogen, carbon dioxide and carbon monoxide, respectively, in the synthesis gas. A module of about 2.0 defines the desired stoichiometric ratio of synthesis gas for the production of methanol. Other important properties of the synthesis gas in methanol production include the carbon monoxide to carbon dioxide ratio and the concentration of inerts in the synthesis gas. A high carbon monoxide to carbon dioxide ratio typically increases the reaction rate of the formation of methanol and the achievable per pass conversion while it concurrently decreases the formation of water thereby reducing the catalyst deactivation rate. A high concentration of inerts in the synthesis gas, such as methane, argon, nitrogen, etc. typically lowers the partial pressure of the active reactants. Since the methanol conversion reaction is exothermic, lower temperatures favor conversion of the synthesis gas to methanol. Pressure will also affect the methanol conversion reaction, with increasing pressure also favoring methanol formation.\nIn many methanol production facilities, the incoming compressed synthesis gas is often mixed with recycled unreacted gas stream to form the synthesis gas stream that is supplied to the methanol converter reactor. A portion of the unreacted gas stream may be purged to prevent the buildup of inerts in the methanol converter reactor. The amount of purge flow typically varies anywhere from 1% to 10% of the total unreacted gas stream and often depends on the amount of inerts in the incoming synthesis gas, with higher level of inerts generally requiring higher purge flows and lower level of inerts generally requiring lower purge flows.\nSome of the prior art uses of the purge stream include use of the hydrogen and/or methane slip in the purge stream as a feed or source of fuel to be used in the front-end steam methane reforming (SMR), partial oxidation (POx), autothermal reforming (ATR) processes. Other prior art has suggested the recovery of hydrogen from the purge stream and mixing the recovered hydrogen with the synthesis gas to improve the module of synthesis gas for methanol production.\nAs used herein, steam methane reforming (SMR) is a catalytic conversion of natural gas, including methane and light hydrocarbons, to synthesis gas containing hydrogen and carbon monoxide by reaction with steam. The reactions are endothermic, requiring significant amount of energy input. The steam methane reforming process is carried out at high temperatures within catalyst filled tubes inside a fired furnace. The amount of steam used is in excess of the reaction stoichiometry requirements, as required to prevent the catalyst from coking. No oxygen is used in steam methane reforming.\nPartial oxidation, on the other hand, is a non-catalytic process where a sub-stoichiometric amount of oxygen is allowed to react with the natural gas creating steam and carbon dioxide at high temperatures. The residual methane is reformed through reactions with the high temperature steam and carbon dioxide to produce synthesis gas. In principle, the partial oxidation reaction can be carried out without any steam addition. Autothermal reforming is a variant of the partial oxidation process, but which uses a catalyst to permit reforming to occur at lower temperatures than the partial oxidation process. Moderate amounts of steam are typically required to prevent the catalyst from coking.\nMany synthesis gas generation methods also employ pre-reforming and secondary reforming. When the feedstock contains significant amounts of heavy hydrocarbons, SMR and ATR processes are typically preceded by a pre-reforming step. As generally known in the art, pre-reforming is a catalyst based process for converting higher hydrocarbons to methane, hydrogen, carbon monoxide and carbon dioxide. The reactions involved in pre-reforming are typically endothermic. Most pre-reformers operate adiabatically, and thus the pre-reformed feedstock typically leaves at a lower temperature than the feedstock entering the pre-reformer. A secondary reforming process conventionally refers to an autothermal reforming process that is fed product from a SMR process. Thus, the teed to a secondary reforming process is primarily synthesis gas from the SMR. Depending on the end application, some natural gas may bypass the SMR process and be directly introduced into the secondary reforming process. Also, when a SMR process is followed by a secondary reforming process, the SMR may operate at a lower temperature, e.g. 650° C. to 800° C. versus 850° C. to 950° C."}
{"text": "The present invention relates to a display device, and more particularly, to a display section having a touch panel built in an organic EL display device.\nTouching the surface of the display device with the user's finger has become common input process for the display device of the smart phone or the tablet. Generally, the touch panel is disposed on the display region. The liquid crystal display panel with built-in touch panel function has been developed for adaption to the liquid crystal display device.\nJapanese Unexamined Patent Application Publication No. 2016-1233 discloses the liquid crystal display panel with built-in touch panel function. The disclosed liquid crystal display panel is configured to dispose one of electrodes of the touch panel at the outer side of the counter substrate, and to allow the common electrode inside the liquid crystal display panel to be used as the other electrode of the touch panel."}
{"text": "In recent years, a secondary battery, which can be charged and discharged, has been widely used as an energy source for wireless mobile devices or an auxiliary power device. In addition, the secondary battery has attracted considerable attention as a power source for electric vehicles (EV), hybrid electric vehicles (HEV), and plug-in hybrid electric vehicles (Plug-in HEV), which have been developed to solve problems, such as air pollution, caused by existing gasoline and diesel vehicles using fossil fuels.\nSmall-sized mobile devices use one or several battery cells for each device. On the other hand, middle or large-sized devices, such as auxiliary power devices or vehicles, use a battery module having a plurality of battery cells electrically connected to each other because high output and large capacity are necessary for the middle or large-sized devices.\nPreferably, the battery module is manufactured so as to have as small a size and weight as possible. For this reason, a prismatic battery or a pouch-shaped battery, which can be stacked with high integration and has a small weight to capacity ratio, is usually used as a battery cell (unit battery) of the middle or large-sized battery module. In particular, much interest is currently focused on the pouch-shaped battery, which uses an aluminum laminate sheet as a sheathing member, because the pouch-shaped battery is lightweight, the manufacturing cost of the pouch-shaped battery is low, and it is easy to modify the shape of the pouch-shaped battery.\nIn general, a high-output, large-capacity battery pack including a plurality of battery modules is configured to have a structure in which a plurality of battery cells or battery modules is stacked while being adjacent to each other, and the battery cell stack or the battery module stack is fixed to a base plate. In this structure, long bolts and additional fixing brackets connected to frames of the battery modules and the base plate in a surface contact fashion are used as means for fixing the battery modules to the base plate.\nIn the fastening structure using the long bolts and the fixing brackets, however, the battery cells may move in a direction in which the battery cells are stacked or a direction perpendicular to the ground when excessive impact is applied to the battery pack with the result that safety of the battery pack is lowered. In addition, the internal space of the battery pack is decreased, and the weight of the battery pack is increased, due to the use of the additional fixing brackets. Furthermore, a process of fastening the fixing brackets increases total process time.\nConsequently, there is a high necessity for a battery module having a novel assembly structure that is capable of solving the above problems."}
{"text": "1. Field of the Invention\nThis invention relates to methods, apparatus and compositions for separating yttrium-90 from strontium-90.\n2. Background of the Related Art\nThe use of radioactive isotopes as diagnostic, imaging and therapeutic agents is a relatively new area of medicine that has flourished in the last fifty years. A number of radioisotopes, primarily beta emitting radionuclides, are finding use in the in vitro treatment of cancers to destroy or sterilize cancer cells. The treatment is administered in a series of cycles to avoid radiotoxicity to other areas of the body, particularly the kidneys and bone marrow. The isotopes of interest are commonly attached to monoclonal antibodies specific for the cancer cells to be treated, thus delivering a dose of radiation directly to a tumor. This technique is termed radioimmunotherapy (RIT) and is increasingly being used to complement existing surgical techniques and chemotherapy.\nIn order to fuel the current research in the use of radionuclides to treat cancers, it is essential that new isotope production methods be developed to increase the availability and decrease the cost of radioisotopes. For medicinal applications, the radioisotope supplied needs to be radiochemically pure to prevent the accidental introduction of unwanted additional radionuclides into a patient, and, preferably, be carrier free. A fundamental aspect of increasing the availability of radioisotopes to medical personnel is the development of new, inexpensive, radiolytically stable materials to allow the necessary separations to be achieved.\n90Y is a high-energy beta emitter that is finding use in the treatment of certain forms of cancer. 90Y decays by pure beta emission, with a half-life (T1/2) of 64 hours, to stable 90Zr. The energetic beta particles (2.3 MeV) can penetrate an average of 0.5 cm in human tissue, with a maximum penetration of up to 1 cm. Consequently, they are useful in the treatment of cancerous tumors like those found in Hodgkin's disease, where tumors are typically between 1 and 5 cm in diameter. The 90Y can be successfully attached to an antibody, which will then transport the 90Y to the targeted tumor.\nIn order to use 90Y in the treatment of cancers, it is necessary to obtain a very pure source of the isotope that is free from the parent 90Sr. This is essential because 90Sr has a 28 year half-life and is likely to accumulate in the bone if inadvertently introduced into the body. The maximum tolerable amount of 90Sr fixed in the bone is only 2 μCi and consequently great care needs to be performed to achieve the necessary Sr/Y separation to ensure minimal introduction of 90Sr into the body during the 90Y radiotherapy.\n90Y is the daughter product of 90Sr, an abundant fission product of 235U, found in nuclear wastes resulting from the reprocessing of spent commercial nuclear fuel and in the separation of 239Pu for weapons manufacture. 90Sr has a half-life of approximately 28 years. The radioactive decay scheme is outlined in Equation 1 below.90Sr(β−)→90Y(β−)→90Zr  (1)\nIn order to obtain a supply of 90Y, it is first necessary to separate 90Sr from other isotopes in the nuclear waste. This can readily be achieved using selective precipitation, ion exchange or solvent extraction techniques to produce a crude 90Sr ‘cow’ for use as a source of 90Y. 90Y can also be produced by the neutron irradiation of 89Y oxide, Y2O3, for a period ranging from several days to a week, but this is expensive and the 90Y product contains large amounts of inactive 89Y, making it unsuitable for medicinal applications.\nThere are a number of methods described in the literature for the separation of the 90Y daughter from the parent 90Sr, including solvent extraction, ion exchange, precipitation and chromatographic procedures. Of these methods, ion exchange techniques have probably received the most attention. However, all of the current methods suffer from drawbacks. For instance, in some separation procedures, the 90Sr is held onto an organic cation exchange resin and the 90Y is eluted using an aqueous complexant solution, such as EDTA, oxalate, lactate, citrate etc. Consequently, the purified 90Y is generated as a complex that is not suitable for the direct labeling of antibodies and requires further processing. Organic ion exchange resins are also prone to radiation damage resulting in a decrease in capacity and the potential release of toxic organic molecules into the 90Y stream as the resin decomposes. Consequently, there is a need for new material and methods to produce pure 90Y.\nThe method disclosed by Bray and Webster in U.S. Pat. No. 5,512,256 uses a solvent extraction process to separate 90Y from 90Sr. A 0.3M solution of di(2-ethylhexyl)phosphoric acid (HDEHP) in n-dodecane is used to extract 90Y from a solution of 90Sr/90Y in 0.3M nitric acid. The HDEHP selectively extracts the 90Y into the organic phase and residual 90Sr can be removed by further washing the organic fraction with fresh 0.3M nitric acid. Although this method is very effective at separating 90Y from 90Sr, multiple steps are required and the recovery of both the 90Sr cow and 90Y fractions requires multiple washing and stripping phases. This produces waste organic and aqueous streams that need to be treated and disposed of safely. There will also be some radiolysis of both the organic complexant and the solvents that will limit their useful life and also may cause the release of unwanted organic species into solution. This is the primary method utilized to produce 90Y in the USA today.\nIn U.S. Pat. No. 5,368,736, Horwitz and Dietz use a multiple step chromatographic process to separate 90Sr from 90Y. The 90Sr stock solution in 3M nitric acid is passed through three strontium selective chromatographic ion exchange columns in series so that the solution exiting the third column contains essentially only 90Y, the 90Sr being retained on the columns. This raw 90Y solution is the passed through a rare earth selective column that selectively extracts the 90Y. The purified 90Y can then be eluted off the column. However, the chromatographic columns contain organic resins that are susceptible to radiation damage and may leach undesirable radiolysis fragments into the purified 90Y stream. Radiation damage is kept to a minimum by loading and then eluting the radioactivity from the columns, but this method also requires the use of a dedicated hot cell facility, necessitating shipment of the purified 90Y to the end user.\nHuntley's U.S. Pat. No. 5,494,647 discloses an ion exchange process for separating 90Y from 90Sr using CHELEX-100® (Bio-Rad Laboratories, Richmond, Calif.), a chelating ion exchange resin. CHELEX-100® is an organic ion exchange resin that consists of iminodiacetic acid groups mounted on a polystyrene/divinyl benzene substrate. The method is designed for use with environmental samples only containing trace amounts of 90Sr, and it is disclosed that the method does not work effectively at high strontium concentrations. The organic resin would also be susceptible to radiation damage and it is doubtful that the method would be able to produce the level of 90Y purity required for medicinal applications.\nTherefore, there is a need for improved methods, apparatus, and compositions for separating yttrium-90 from strontium-90. It would be desirable if the compositions were highly radiation resistant, thermally stable, chemically stable, and non-toxic. It would be even more desirable if the compositions and methods provided very high affinities for strontium-90 and very low affinities for yttrium-90."}
{"text": "Non-patent document 1 discloses a conventional inverter modulation method for driving an electric locomotive. In the inverter modulation method disclosed in Non-patent document 1 multi-pulse asynchronous modulation is used when inverter frequency is low, three-pulse synchronous modulation is used when the inverter frequency increases to a predetermined level, and single-pulse synchronous modulation is used when or after the output voltage reaches a termination voltage (rated output: 100%).\nNon-patent document 1:“Instruction to inverter-controlled electric locomotive” Hideki Iida, Atsushi Kaga, Denkisyakenkyukai"}
{"text": "Modern computer servers can have terabytes of physical memory. When a system panics and there is a desire to collect diagnostic information, the process of dumping physical memory to a file system can be very time intensive. The resulting downtime can have significant impact on a business's bottom line. Consequently, some server owners eschew fault diagnosis in favor of shorter downtime, a decision that may lead to more costly system failure down the road."}
{"text": "There has been known a radio communication system configured in assuming that a base station performs radio communication with a mobile terminal (mobile station) included in a device that travels on a predetermined route, such as a train and an automobile. As a method of configurating such a radio communication system, there has been proposed a method of configuring a communicable range (cell) in an elongated shape along a traveling route with a plurality of directional beams using a distributed antenna base station to which a plurality of directional antennas, which is distributed along the traveling route of a mobile device (mobile station), is connected (Non Patent Literature 1). When the shape of the cell is formed to be elongated along the traveling route, a cell length can be increased while interference with another radio communication system is suppressed. Such a distributed antenna system can be implemented in a form in which a plurality of antennas is arranged from one communication modem using a transmission medium such as an optical fiber.\nMeanwhile, there has been known a method of measuring (ranging) a distance between a base station and a mobile station as a method of detecting a location of a mobile device, such as detection of a traveling location of a train. As such a method of measuring a distance, there has been known a method of measuring a distance from the time required for communication between a radio device of the distance measuring side (e.g., base station) and a radio device of the distance measured side (e.g., mobile station) (Non Patent Literature 2)."}
{"text": "A simulation is an imitation of some real thing or process. Simulation can be used to show the unfolding and eventual real effects of alternative conditions and courses of action or inaction. Simulations are used in many contexts, including the modeling of natural or human systems, in order to gain (or provide) insight into their functioning. Contexts in which simulation is used include performance optimization of a wide range of technologies, safety engineering, testing, training and education.\nA computer simulation is an attempt to model a real-life or hypothetical situation on a computer so that, among others, it can be studied to see how the system works. By changing variables, predictions may be made about the behavior of the system. Computer simulation is often used as an adjunct to, or substitution for, modeling systems for which simple closed form analytic solutions are not possible.\nSimulation, particularly computer simulation, is often used in training. This includes the training of civilians as well as military personnel. This often occurs when it is prohibitively expensive or too dangerous to allow trainees to use real equipment in the real world. In such situations they spend time learning valuable lessons in a “safe” virtual environment. Often the convenience is to permit exploration, mistakes, and exposure to rare events during training for a safety-critical system.\nFor example, a flight simulator is used to train pilots on the ground. It permits a pilot to “crash” a simulated “aircraft” without harming the instructor, pilot and aircraft. Flight simulators are often used to train pilots to operate aircraft in extremely hazardous situations, such as landings with no engines, or complete electrical or hydraulic failures. The most advanced simulators have high-fidelity visual systems and hydraulic motion systems. The simulator is normally cheaper to operate than a real aircraft (fuel, wear and tear, removing aircraft from revenue generating service).\nBearing resemblance to flight simulators, marine simulators train a ship's personnel. Simulators like these are mostly used to simulate large or complex vessels, such as cruise ships, dredging ships, or submarines. They often consist of a replication of a ships' bridge, with operating console(s), and a number of screens on which the virtual surroundings are projected.\nStrategy games—both traditional and modern—may be viewed as simulations of abstracted decision-making for the purpose of training military and political leaders. Many video games are also simulators.\nMedical simulators are increasingly being developed and deployed to teach therapeutic and diagnostic procedures as well as medical concepts and decision making to personnel in the health professions. Simulators have been developed for training procedures ranging from the basics such as blood draw, to laparoscopic surgery and trauma care.\nSome medical simulations are developed to be widely distributed (such as web-enabled simulations that can be viewed via standard web browsers) and can be interacted with using standard computer interfaces, such as the keyboard and mouse.\nComputer simulations have the advantage of allowing a student to make judgments, and also to make errors and learn from these errors. The process of iterative learning through assessment, evaluation, decision making, and error correction creates a much stronger learning environment than passive instruction.\nThus, realistic interactive simulations can be an important tool for learning difficult, dangerous, rare or not currently accessible information. Simulations can be performed to learn a task, play a game, and even to take a tour of a particular location. One problem encountered in implementing display-based or screen based simulations is the ability to provide the desired information within the finite space of a display such as a computer screen. Some approaches involve using a coarse image resolution in order to provide an extensive environment on one screen. Another approach, when finer resolution is required and only a smaller portion of the environment can be displayed on the screen at a given time, involves splitting a representation of an entire environment into multiple graphical representations. Navigation between various windows of an environment that has been split into multiple graphical representations is generally required for simulating an extensive functional environment. Unfortunately, splitting a representation of an entire environment into multiple graphical representations may introduce the need for pull-down menus and navigation between the various windows that (a) reduce user friendliness of the graphical user interface, (b) interfere with the flow of the simulation, (c) interrupt the suspension of disbelief necessary for a user to get immersed in a simulation, and (d) degrade the user's perception of the spatial relationship between the different component views of the entire simulated system and interfere with situational awareness.\nAlternatively, a collage or patchwork of graphics representing different parts of an environment has been crammed onto a single screen. Sometimes these disparate graphics may confusingly be at different scales or taken from different perspectives and viewing angles and rendered in different formats, potentially making it difficult for users to appreciate the position, size and appearance of the components relative to each other.\nAccordingly, there is a need for a dynamic, interactive simulation system having excellent resolution and yet extensive scope and minimally interrupted suspension of disbelief."}
{"text": "The present invention relates to voltage regulators (VRs), and more specifically to placement of VR components in a system.\nVoltage regulators are used in systems such as a personal computer (PC) (e.g., a desktop computer, server computer, notebook computer and the like) to receive input direct current (DC) voltages of a given voltage and convert and regulate such DC voltages to one or more regulated voltage levels required by various system components, such as integrated circuits (ICs) and the like.\nIn a typical system, for example, a desktop PC, a motherboard is used to support various system components, including ICs, connectors, VR components, and the like. Such VR components may include output inductors, bulk capacitors, metal oxide silicon field effect transistors (MOSFETs), driver ICs, and the like. Typically, the VR components are placed on a primary side (e.g., a topside) of the motherboard. These VR components must be placed outside an IC device's keepout zone (i.e., the IC's footprint, including any socket or heatsink retention). This placement can result in larger loadline lengths and therefore higher loadline impedance (i.e., capacitance, inductance, and resistance). In such manner, VR components may be located many centimeters away from an intended load (e.g., an IC). Accordingly, motherboard/package lateral travel dominates, increasing loadline length and therefore impedance.\nInstead of the above described placement of VR components, other systems use some type of additional VR daughter module (i.e., a separate circuit board) that is plugged into the motherboard or an IC device. However, such additional circuit boards increase cost and complexity, and further increase the size of a given form factor. Furthermore, such designs typically provide inferior performance. A need thus exists to provide VR components that have reduced loadline length and impedance."}
{"text": "Computer microchips permeate almost all aspects of modern life. From the single microchips embedded into our pets to the myriad of microchips “flying” our aircraft and keeping our cars running, our reliance on semiconductor device technology cannot be minimized. A topic of great importance, considering all of the critical operations controlled by semiconductor devices, is the durability and reliability of those devices. One issue that has been at the forefront of semiconductor device design is protecting devices from electrostatic discharge (ESD). ESD is defined as the transfer of charge between bodies at different electrical potentials. ESD is a serious issue in solid state electronics. Integrated circuits are made from semiconductor materials such as silicon and insulating materials such as silicon dioxide. When subjected to high voltages, many semiconductor materials can suffer permanent damage, changing the electrical characteristics, degrading, or destroying it. It may also upset the normal operation of an electronic system, causing equipment malfunction or failure.\nDamage to electronic devices from ESD can occur at any point from manufacture to field service. Damage may result from improperly handling the devices in uncontrolled surroundings or when substandard ESD control practices are used. Generally, ESD damage is classified as either a catastrophic failure or a latent defect.\nA catastrophic failure occurs when an electronic device no longer functions after exposure to an ESD event. An ESD event may cause a metal melt, junction breakdown, or even oxide failure. A device's circuitry can, therefore, be permanently damaged causing the device to fail. Such failures can usually be detected when a device is tested before shipment. However, if the ESD event occurs after testing, the damage may go undetected until the device fails in operation.\nA latent defect, on the other hand, is typically more difficult to identify. A device that is exposed to an ESD event may be partially degraded, yet continue to perform its intended function. However, this degradation generally reduces the operating life of the device dramatically. A product or system that incorporates a device with latent defects may experience premature failure after the user places the device into service. Such failures are usually costly to repair and in some applications may create personnel hazards.\nVarious external solutions and procedures have been developed for preventing ESD damage during fabrication and device manufacturing. Manufacturers often implement electrostatic protective areas (EPAs). An EPA can be a small working station or a large manufacturing area. The main principles of an EPA are: (1) there are no highly charging materials in the vicinity of ESD sensitive electronics; (2) all conductive materials are grounded; and (3) workers are grounded. Adherence to these principles may prevent charge build-up on ESD sensitive electronics. International standards are used to define typical EPA and can be found, for example, from International Electrotechnical Commission (IEC) or American National Standards Institute (ANSI).\nESD prevention within an EPA may include: using appropriate ESD-safe packing material, using conductive filaments on garments worn by assembly workers, using conducting wrist straps and foot-straps to prevent high voltages from accumulating on workers' bodies, using anti-static mats or conductive flooring materials to conduct harmful electric charges away from the work area, and using humidity control. Humid conditions prevent electrostatic charge generation because the thin layer of moisture that accumulates on most surfaces serves to dissipate electric charges. Ion generators are also sometimes used to inject ions into the ambient airstream. Ionization systems help to neutralize charged surface regions on insulating or dielectric materials.\nIn addition to such external ESD prevention mechanisms, chip designers have also incorporated ESD protection internally into the design of the device. Various methods and configurations for adding N+ or P+ doped regions have been implemented in field effect transistor (FET) devices in association with source/drain (S/D) and gate regions. One common method is to provide a zero-space N+ S/D implant region between contacts.\nA disadvantage of the zero-space N+ S/D implant region and similar configurations is that they do not typically provide good and qualified ESD protection for high voltage applications."}
{"text": "Solid state hybrid switches have been known heretofore. For example, C. G. Chen et al U.S. Pat. No. 4,420,784, dated Dec. 13, 1983, shows one form of hybrid D.C. power controller of the relay/circuit breaker type that uses a hybrid arrangement of hard contacts and power FETs in cooperative functional combination which is especially adapted for switching 270 volt D.C. power systems in the low atmospheric pressure environments. Also, A. R. Perrins U.S. Pat. No. 3,330,992, dated July 11, 1967, shows an A.C. load energizing circuit having a pair of reverse-parallel controlled rectifiers, a first switch for applying anode voltage to these controlled rectifiers, a second switch for gating these rectifiers into conduction to establish the circuit to the load and a third switch for shunting the controlled rectifiers. A variety of zero voltage solid state switching circuits have also been known. However, these prior switching circuits have been handicapped by certain disadvantages such as, for example, being rather complex in structure with large numbers of components, insufficient contact life, current leakage in the off condition or inadequate switching performance. Accordingly, it has been found desirable to provide an improved solid state hybrid switch that overcomes the aforesaid advantages."}
{"text": "Inventory control for long-term storage items has become an increasingly important field of endeavor. One particularly noteworthy aspect of inventory control is environmental control of the storage location, whether the commodities are stored in a depot, warehouse or retail establishment. One parameter that should be monitored during long-term storage is humidity, because such changes in humidity can have drastic and long-reaching impact on many different commodities. If the proper humidity range for the stored object is not maintained during storage, the useful lifetime, capability, safety, and other characteristics of the stored object can be deleteriously degraded or lost, which can lead to a number of undesirable, and possibly dangerous, results.\nFor monitoring a long-term storage condition, it is a disadvantage for any humidity monitor to require a battery power source, because any battery has a finite useful life over a certain period of time. Therefore, it is essential that a long-term humidity monitor must operate without battery power.\nOne type of prior art humidity monitor that does not require electrical power is the dial hygrometer. FIG. 1 is a perspective view of a prior art dial hygrometer with a needle that moves in response to changes in humidity, which is set against a clock-type face. While the prior art dial hygrometer is useful for indicating current humidity through movement of the needle, it does not record the history of humidity changes during storage. The inability of the prior art dial hygrometer to record the history of humidity changes during storage is a key disadvantage of this device. Thus, there has been a long-felt need for a real-time humidity monitor that can also record humidity change history.\nRadio Frequency Identification (RFID) is an emerging area of technology with numerous applications, particularly in the field of inventory control. RFID is an automatic identification method, relying on storing and remotely retrieving data using devices that are called RFID tags or transponders. An RFID tag is a small object that can be attached to, or incorporated into, a product, animal, person or other stored object and contains silicon chips and antennas to enable it to receive and respond to radio-frequency queries from an RFID transceiver. When an RFID tag is combined with one or more sensors, the RFID tag can report the information that is measured by the sensors. A passive RFID tag requires no internal power source, while active RFID tags do require a power source. A passive RFID tag is better for inventory control for long-term storage items because the use of battery power to operate the RFID tag is ordinarily unacceptable. In the RFID arrangement, the device receives electrical power from the RF when it is being read, therefore for such applications, a passive RFID tag combined with a sensor that also does not need power for its operation is critical. However, there is no currently available RFID tag that can effectively monitor changes in humidity. Thus there has been a long-felt need for a humidity monitor using RFID technology that not only responds to changes in humidity but can also provide a real-time record of those humidity changes. Up until now, the long-felt need for an all-humidity-sensitive RFID tag has not been met."}
{"text": "1. Reference to Appendix\nAppendix A is a part of the present disclosure, and is incorporated herein by reference in its entirety.\nA portion of the disclosure of this patent document contains material which is subject to copyright protection. The copyright owner has no objection to the facsimile reproduction by anyone of the patent document or the patent disclosure, as it appears in the Patent and Trademark Office patent files or records, but otherwise reserves all copyright rights whatsoever.\n2. Field of the Invention\nThe present invention is related generally to monitoring operation of an integrated circuit and in particular to an integrated circuit output port for transporting information concerning any user selected operation of the integrated circuit from the integrated circuit.\n3. Description of Related Art\nAs the form factor of hard disk drives become smaller and as the level integration for all disk drives becomes greater, the number of functions performed by a single integrated circuit in a disk drive increases substantially. As the number of functions performed by an integrated circuit increases, i.e., as the integration level increases, many of the functions are performed completely internally to the integrated circuit.\nFor functions that are performed completely internally, neither input, output, nor intermediate signals are available at pins of the integrated circuit. Further, the integrated circuit package size limits the number of available pins. The limited number of pins makes it unfeasible to route representative signals for each of the internal functions to a pin. Consequently, monitoring and testing of highly integrated circuits becomes problematic.\nIn the area of disk drive controllers, integration is becoming increasingly greater. Disk drive controller circuits that previously occupied a plug-in board are now packaged in a single integrated circuit. The plug-in boards had space for test points, switches for generating test signals, and other commonly known features that permitted relatively easy verification of the operation of the disk controller. When the same functionality is provided in an integrated circuit, there is no space for test points, switches and similar circuitry.\nTherefore, typically, computer simulations are used to verify the operation of circuitry that is totally internal to an integrated circuit. Unfortunately, such simulations do not provide reliable information on the actual application performance of the integrated circuit. To obtain reliable application performance data, a way is needed to obtain information concerning all of the circuits within the integrated circuit while the circuits are in operation."}
{"text": "The present invention is directed to a cable sleeve composed of a longitudinally-divided cable housing having a sealing system in the parting plane of the housing and having cable introduction areas in the end faces of the housing.\nU.S. Pat. No. 4,117,259, whose disclosure is incorporated herein by reference thereto and which claims priority from German Application 25 15 939, discloses a cable sleeve that is composed of two half-shells, wherein the cable introduction openings lie in a parting plane of the sleeve housing. Given such a sleeve form, the cable introduction openings must be completely renewed after the sleeve has been opened."}
{"text": "Commercial aircraft operators and passengers are increasingly requesting in-flight entertainment (IFE) and electrical power outlets to enable passengers and/or crew members to operate various personal electronic devices, such as laptop computers, at their seats. To provide passengers and/or crew members with such features, electrical power and data communications capability must be delivered to each seat. Power and data are delivered to aircraft passenger seats via one or more cables that extend from a power or data source to the passenger seats. The cables typically extend up through a floor access panel and along one or more seat tracks recessed within the floor of the aircraft passenger cabin. The seat track is covered by a seat track cover.\nThe cover is typically secured directly to the seat track or to the portions of the aircraft passenger cabin floor that surround the seat track. However, in some applications the seat track is planar and does not provide a surface that will permit direct engagement with the seat track cover. Therefore, there is a need for a system for securing seat track covers to seat tracks with planar upper surfaces."}
{"text": "1. Field of the Invention\nThe invention relates to novel lead structures connecting to die that provide for electromagnetic interference reductions. Crosstalk between leads, as well as susceptibility to noise created by electromagnetic emissions internal or external to the package, are reduced through the implementation of the invention.\nFurther, the invention relates to novel multiple ground planes formed to connect dielectric coated leads connecting to one or more dies.\n2. Description of Related Art\nElectromagnetic interference leading to reduced performance is an increasingly common problem for packaged die, particularly for die having input/output (IO) operating at Gigahertz frequencies. Many integrated circuits generate undesirable amounts of EMI. Typically, the noise generated by the integrated circuit originates from the die and its connections to the pins through the package. As the EMI is coupled to neighboring components and integrated circuits, it interferes with their individual performance which may, in turn, affect the overall performance of a system. Because of the negative effects of EMI and because the level of acceptable radiated EMI is subject to strict regulatory limits, it is desirable to contain or suppress the EMI generated by an integrated circuit.\nSolutions such as separation of leads or isolation with shields are not always available or sufficient. Furthermore, EMI solutions at the IC package level are often ignored because the main concerns at that level are signal integrity and functionality. It would be beneficial to have an EMI solution at the package level because it would help reduce the need for “downstream” or add-on solutions."}
{"text": "1. Field of the Invention\nThe present invention relates to a for an ileostomy bag, a colostomy bag or the like. This bag includes a deodorizing means which is self contained and easily dispersed into the bag to enable reuse.\nThis invention relates to a colostomy bag, more particularly, to a bag that is adapted to house body fluids and waste matter and thereby dispersing deodorizer into said bag.\n2. Description of the Prior Art\nIn the treatment of human diseases and ailments, it is sometimes necessary to form an opening in the patient's anatomy and to maintain that opening for an indefinite period of time. For instance, diseases involving different parts of the gastro-intestinal and urinary tract can result in a patient being left with an abdominal stoma. The three most common types of abdominal stomas are the colostomy, the ileostomy and the ileal conduit. In the case of an ileostomy, ileal conduit and many colostomy operations, the patient is unable to control the passage of bodily waste material and must rely upon an appliance attached to the body to collect this material.\nVarious appliances have been proposed in the prior art. A majority of these appliances involve the use of deodorant gels, tablets and liquids or other deodorizing materials which must be placed within the ostomy bag. However, these materials are often messy to apply and are often inadequate to prevent odors. For instance, if a user perspires, such as after exercise or on a hot day, these deodorizing materials are often ineffective in providing a complete deodorizing effect.\nFurther, many of the prior art arrangements are rather complicated to operate because they require a user to assemble various components or open the ostomy bag and insert the deodorizer. Thus, the user tends to avoid changing the ileostomy or colostomy bag and bacterial infections or other problems may arise. Further, as a complicated sealing arrangement is required, it is often difficult for the bag to be properly attached to the multi-piece connector and leaks may therefore arise.\nFurthermore, such complicated deodorizing arrangements may be time consuming and may be difficult for certain users, such as elderly, handicapped or impaired persons or those with arthritis, to use. Over time, prior art constructions are subject to wear. Due to their constructions, this wear can result in leaking of the prior art connectors. Furthermore, the manufacturing cost for these prior art devices tend to rise because of the need for manufacturing and packaging of several different components.\nAccordingly, a need exists in the art for a simple and effective bag deodorizing system which is self-contained and easy to operate, which avoids the opening and addition of deodorants without having to carry said deodorants on oneself for application within said ostomy bag.\nThe prior art teaches a variety of colostomy appliances, for example, as disclosed in U.S. Pat. Nos.: 2,546,779; 2,679,248; 2,684,675; 2,814,295; 2,896,625; 3,125,093; 3,695,268; 3,805,789; 3,826,262; and others. None of the foregoing, nor the prior art described below however, provide for a simple yet effective means of dispersing deodorant into a ostomy pouch for reuse of same hence, eliminating odor."}
{"text": "This invention relates generally to an ultrasonic gas measuring device, for detecting and determining the concentration of a gas being measured, and incorporates means for providing a more effective and efficient determination of the concentration of any gas being detected whether said gas is in the process of being manufactured, or utilized, such as in medical, industrial, or other applications.\nVarious gas measuring means for determining the concentration of a gas within a particular specimen, such as the concentration of oxygen as used in medical applications, have been available in the art. For example, a Douglas Scientific Products, Inc. has been marketing an oxygen concentrator monitoring device, identified as the \"Oxyalert,\" and which feeds a sample of the oxygen being measured into one area of a transducer assembly, and its sample chamber, and then at a remote location has an exit point for the measured gas, after passing it through an ultrasonic detector that senses the gas composition, and the concentration of specific gaseous components, within its device.\nHudson Oxygen Therapy Sales Co. markets an oxygen concentrator monitoring and regulation assembly, as shown in its U.S. Pat. No. 4,516,424, wherein the concentration of oxygen is detected by passing a gas over a fuel cell terminal that analyzes the gas being measured, by generating a signal as a result of exposure of the gas to the cell terminal, as a means for providing a determination of the concentration of the gas exposed thereto. Other patents to Hudson pertaining to the same technology are shown in the U.S. Pat. No. 4,627,860, U.S. Pat. No. 4,516,424, U.S. Pat. No. 4,561,287, and the U.S. Pat. No. 4,648,888.\nThe United States patent to Kniazuk, et al, U.S. Pat. No. 2,984,097, discloses a gas measuring apparatus, in which embodiment the development was used to determine the consumption of oxygen being made by a test specimen, and then measuring that degree of oxygen consumed by passing it through what is identified as a sonic generator.\nThe patent to Mechlenburg, U.S. Pat. No. 4,616,501, discloses a system that provides an indication of the concentration of at least one gas within a chamber, wherein a pair of operating frequencies generated from mechanical radiant energy are detected, their amplitudes measured, and then compared to a modeled amplitude response representative of a plurality of that gas's concentration to determine the specific concentration of the sample being tested.\nThe patent to Noguchi, et al, U.S. Pat. No. 4,662,212, defines a measuring instrument for concentration of gas, that does utilize and ultrasonic vibrator, but wherein the ultrasonic sensor is made humidity proof by utilizing various sealing materials on its outer surface by depositing a film formation by vacuum evaporation through the usage of select elements or compounds for achieving its structure. The device, though, is used for measuring the concentration of a gas passing through the instrument.\nThe U.K. published patent application No. GB2,087,559A, to Baird, et al, discloses a means for determining gas compositions accoustically, such as by an ultrasonic isolator, for determining and detecting the measurements of a fluid, which presumably includes gases, as well as liquids.\nThe patent No. SU112270A, of the Soviet Union, discloses a two/phase substance gas concentration acoustic measuring unit, for use for testing a substance, including a gaseous concentration. Phase-sensitive characteristics of the unit are determined during the change of the gas concentration, with the device being utilized upon a pipeline for measuring the gaseous substance and the concentration of that gas.\nThe patent to Ringwall, et al, U.S. Pat. No. 3,805,590, discloses an oxygen partial pressure sensor, which utilizes the detection of pressure for determining the amount of oxygen in a helium-oxygen gaseous mixture, and does define the usage of acoustic signal phase shifting to generate electrical signals by which the apparatus maintains oxygen partial pressure within required limits in its delivery of the gas apparently for breathing purposes at various water depths.\nOther United States patents relating to detection of gas quantities include the United States patent to Advani, et al, U.S. Pat. No. 4,462,246; the United States patent to Terhune, U.S. Pat. No. 4,520,654; and, the United States patent to Ogura, et al, U.S. Pat. No. 4,581,942, relating to a measuring conduit for flow rate and concentration of fluid.\nIn addition to the foregoing, a company by the name of Ceramatec, Inc., of Salt Lake City, Utah, has a galvanic oxygen sensor device, which utilizes the fuel cell means for detection in determining the percentage of oxygen within a gaseous mixture.\nThe current invention, as distinct from the apparatuses and developments described in the prior art devices, provides for a more precise analysis of the concentration of a gas or gases within a gaseous mixture, primarily through the regulation of the gas to be measured as it passes into and through the sensing means, such as by achieving its laminar or divided flow, and through the incorporation of temperature, humidity, and pressure compensation means, furnishes a far more consistent and accurate analysis of the specific concentration of any one of the gases being detected, particularly oxygen, than can be achieved through the usage of the individual or combined prior art teachings."}
{"text": "There are many electrical connectors which are known from the published prior art or the marketplace. These connectors seek to connect together electrical conductors without soldering and often without the use of tools. Connectors exist for multistranded insulated wires or cables as well as coaxial cables.\nSeveral such connectors are sold by Swenco Products, Inc. under the mark POSI-LOCK®. Many of these connectors are illustrated in U.S. Pat. Nos. 5,228,875; 5,868,589; 6,358,103 B1; 6,494,753 B1; 6,568,952 B1; 6,692,313 B1; 6,695,653 B1; 6,814,630 B1; 6,830,491 B1; 6,851,966 B1; 6,866,550 B1; and U.S. Patent Application Pub. No. US2004/0192121 A1. These connectors usually require stripping the insulation off of a terminal portion of the wire, and all are connected together by twisting a cap onto a connector body. But helical twisting motions of a multistranded conductor as it is being connected often torsionally stress the metallic strands sought to be connected, resulting in a less than optimum physical and electrical connection. A need therefore persists for connectors which can make a secure electrical connection to a multistranded insulated electrical conductor without twisting one part onto another."}
{"text": "It is quite known already to prepare acetylene polymers by using a catalytic system consisting of a compound of a transition metal and of an organometal compound of the elements of groups 1A, 2A, and 3A of the Periodic System.\nA preparation of this type is found described for the first time by G. NATTA et al. in 'Rend. Acc. Naz. Lincei, Scienze Fis. Mat. Nat.: 25, 3 (1958), according to which crystalline acetylene polymers are obtained in the form of a powder by using catalysts based on titanium and aluminum trialkyl alcoholates.\nMore particularly there is known the possibility of obtaining the polyacetylene directly in the polymerization in the form of a film (T. Iyo, H. Shirakawa & S. Ikeda, J. Polym. Sci. Polym. Chem., Ed. 13, 11 (1974).\nThis method is preferred in general even if with the same it is not possible to obtain films with good mechanical properties.\nIn the published Italian Patent Appl. No. 22722 A/82 there are described new polyacetylene films endowed with particular characteristics, such as the elongation at break, the breaking load and the molecular orientation under stretch. Such films are prepared by using as catalysts special titanium compounds, mainly dialcoxy-dicresoxy titanium in combination with aluminum trialkyls. In general, the values of apparent density of the polyacetylene prepared according to the known methods of the Prior Art prove to be comprised between 0.4 and 0.5 g/cc, with values of the surface area greater than 70 m.sup.2 /g.\nA considerable drawback of the polyacetylenes prepared so far is represented by their instability to oxidation, which considerably reduces the possibility of their being used, for instance, in the manufacture of electrodes for primary and secondary batteries."}
{"text": "The invention relates to a device for converting series of input data elements to series of output data elements, the device being provided with a memory means for containing the series of input and output data elements.\nSuch a device is known from paragraph 10.10 in the book entitled xe2x80x9cTheory and application of digital signal processingxe2x80x9d by the authors Lawrence R. Rabiner and Bernard Gold. In digital signal processing the discrete Fourier transform (DFT) plays an important part. By means of this transform, a description of a signal in the time domain can be converted to a description of the same signal in the frequency domain. This has the advantage that subsequently specific signal processing operations, which can only be performed in a very complex manner in the time domain, can be performed in a relatively simple manner.\nSince performing a DFT requires very many arithmetical processes, through the years many devices have been designed by means of which a DFT can be calculated rapidly and efficiently. Often these devices are provided with a number of so-called butterflies, which are capable of rapidly performing specific sub-processes necessary for the DFT. In such a butterfly, a number of input data elements are converted to an equally large number of output data elements by doing additions and multiplications, said multiplications being performed by means of a multiplying factor. These multiplication factors are alternatively referred to as twiddle factors or WNk.\nThe device mentioned in said book is provided with a RAM-memory, a butterfly and a permutation network. The RAM-memory is embodied so that each memory word can contain a series of data elements, the number of data elements in the series being equal to the number of data elements which can be converted by the butterfly. If, for example, in the known device a radix-4 butterfly is employed, that is a butterfly capable of converting four input data elements to four output data elements, then the RAM-memory is organized so that each memory word can contain a series of four data elements.\nFrom the RAM-memory, a series of all necessary input data elements is supplied to the butterfly by reading a memory word. In the butterfly, these input data elements are converted to output data elements. The multiplication factors necessary for this purpose are read from a ROM-memory. Subsequently, in the permutation network, these output data elements and other output data elements computed by the butterfly are distributed over a number of series of output data elements. These series of output data elements are finally written in the RAM-memory and can be used at a later stage as series of input data elements. By repeatedly performing a radix-B butterfly, in which the output data elements of the butterfly serve in a different order as input data elements for the butterfly, an N-point DFT can be calculated, B being smaller than N. In this manner, for example, a 32-point DFT can be calculated by means of a radix-4 butterfly.\nThe known device is relatively slow.\nIt is an object of the invention to provide a device of the type mentioned in the opening paragraph, which is capable of performing the operation of converting the series of input data elements to series of output data elements more rapidly.\nTo achieve this, the device in accordance with the invention is characterized in that the memory means is embodied so as to read a series of input data elements and write a series of output data elements during a single clock cycle. During the calculation of an N-point DFT, frequently series of input data elements must be read from the memory means and series of output data elements must be written in the memory means. Generally, this takes up one clock cycle for each read or write operation. By embodying the memory means so that during a single clock cycle both a series of input data elements can be read and a series of output data elements can be written, for example by providing the memory means with a so-called read-modify-write function, the access to the memory means is accelerated. This has the additional advantage that the device consumes less energy.\nAn embodiment of the device in accordance with the invention is characterized in that an order in which the series of input data elements are converted to series of output data elements is selected so that multiplication factors required in the butterfly remain constant for as long as possible. If these multiplication factors are stored in a memory, it is achieved by means of this measure that the device has to read a multiplication factor from the memory only a minimum number of times. If, however, the multiplication factors have to be computed by the device, it is achieved by this measure that the device has to compute a multiplication factor only a minimum number of times. In either case, this may lead to a reduced energy consumption by the device and an increase of the speed with which the device converts input data elements to output data elements.\nA further embodiment of the device in accordance with the invention is characterized in that the device comprises at least a Viterbi butterfly. The way in which a DFT is computed has many similarities to the way in which a so-called Viterbi decoding algorithm is computed. By using at least one Viterbi butterfly, the device in accordance with the invention can also suitably be used to realize a Viterbi decoder.\nAnother embodiment of the device in accordance with the invention is characterized in that an order in which the series of input data elements are converted to series of output data elements is chosen so that writing a series of output data elements after reading a necessary series of input data elements occurs as much as possible. By choosing the order in which the series of input data elements are converted to series of output data elements to be such that reading a series of input data elements and writing a series of output data elements during a single clock cycle occurs as much as possible, it is achieved that the memory means must be accessed a minimum number of times. This measure too can lead to a reduction of the energy consumption by the device and an increase of the speed with which the device converts input data elements to output data elements.\nThese and other aspects of the invention will be apparent from and elucidated with reference to the embodiments described hereinafter."}
{"text": "Normal growth, as well as tissue repair and remodeling, require specific and delicate control of activating growth factors and their receptors. Fibroblast Growth Factors (FGFs) constitute a family of over twenty structurally related polypeptides that are developmentally regulated and expressed in a wide variety of tissues. FGFs stimulate proliferation, cell migration and differentiation and play a major role in skeletal and limb development, wound healing, tissue repair, hematopoiesis, angiogenesis, and tumorigenesis (reviewed in Ornitz, Novartis Found Symp 232: 63-76; discussion 76-80, 272-82 (2001)).\nThe biological action of FGFs is mediated by specific cell surface receptors belonging to the Receptor Protein Tyrosine Kinase (RPTK) family of protein kinases. These proteins consist of an extracellular ligand binding domain, a single transmembrane domain and an intracellular tyrosine kinase domain which undergoes phosphorylation upon binding of FGF. Four FGFRs have been identified to date: FGFR1 (also called Flg, fms-like gene, flt-2, bFGFR, N-bFGFR or Cek1), FGFR2 (also called Bek-Bacterial Expressed Kinase-, KGFR, Ksam, KsamI and Cek3), FGFR3 (also called Cek2) and FGFR4. All mature FGFRs share a common structure consisting of an amino terminal signal peptide, three extracellular immunoglobulin-like domains (Ig domain I, Ig domain II, Ig domain III), with an acidic region between Ig domains (the “acidic box” domain), a transmembrane domain, and intracellular kinase domains (Ullrich and Schlessinger, Cell 61: 203, 1990; Johnson and Williams (1992) Adv. Cancer Res. 60: 1-41). The distinct FGFR isoforms have different binding affinities for the different FGF ligands.\nAlterations in FGFRs have been associated with a number of human cancers including myeloma, breast, stomach, colon, bladder, pancreatic and hepatocellular carcinomas. Recently, it was reported that FGFR4 may play an important role in liver cancer in particular (PLoS One, 2012, volume 7, 36713). Other studies have also implicated FGFR4 or its ligand FGF19 in other cancer types including breast, glioblastoma, prostate, rhabdomyosarcoma, gastric, ovarian, lung, colon (Int. J. Cancer 1993; 54:378-382; Oncogene 2010; 29:1543-1552; Cancer Res 2010; 70:802-812; Cancer Res 2011; 71:4550-4561; Clin Cancer Res 2004; 10:6169-6178; Cancer Res 2013; 73:2551-2562; Clin Cancer Res 2012; 18:3780-3790; J. Clin. Invest. 2009; 119:3395-3407; Ann Surg Oncol 2010; 17:3354-61; Cancer 2011; 117:5304-13; Clin Cancer Res 2013; 19:809-820; PNAS 2013; 110:12426-12431; Oncogene 2008; 27:85-97).\nTherapies involving FGFR4 blocking antibodies have been described for instance in WO2009/009173, WO2007/136893, WO2012/138975, WO2010/026291, WO2008/052798 and WO2010/004204. WO2014/144737 and WO2014/011900 also describe low molecular weight FGFR4 inhibitors."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor process, and more particularly to a method of fabricating a semiconductor device and a method of adjusting a lattice distance in the device channel region.\n2. Description of the Related Art\nIn the early days, a metal-oxide-semiconductor (MOS) device is composed of a metal gate electrode, a gate dielectric layer, and a semiconductor substrate. Because the adhesion of most metals to silicon is unsatisfactory, currently, the material of the gate electrode is polysilicon. The application of a polysilicon gate electrode, however, incurs other issues. For example, the device performance decays due to the high resistance of the polysilicon. Accordingly, with the present technology, after forming the device, a salicide process is performed to form metal silicide on the gate electrode and the source/drain regions to reduce the resistance of the device.\nIn another aspect, a chip usually comprises a device area and a peripheral circuit area, wherein, devices in the device area include, for example, memory devices, and electro-static discharge (ESD) protection circuits. Devices in the peripheral circuit area comprise, for example, logic devices. The devices in the device area require high resistances than those in the peripheral circuit area. During the salicide process above, a block layer is used to cover the area on which the metal silicide is not going to be formed. Because the area covered by the block layer does not require additional film layer to prevent the formation of metal silicide, the block layer is also called a self-aligned salicide block layer (SAB) layer.\nFIGS. 1A-1E are schematic cross sectional views showing the progression of a prior art method of fabricating a semiconductor device. Referring to FIG. 1A, a substrate 100 is provided. The substrate 100 includes a device area 102, and a peripheral circuit area 104. Gate structures 106 and 108 are formed over the substrate 100 of the device area 102 and the peripheral circuit area 104, respectively. Lightly-doped regions 110 and 112 are formed in the substrate 100 and adjacent to the sidewalls of the gate structures 106 and 108.\nReferring to FIG. 1B, after forming spacers 114 on the sidewalls of the gate structures 106 and 108, the source regions 116a and 118a, and the drain regions 116b and 118b are formed in the substrate 100 and adjacent to the spacers 114. An anneal process 120 is performed to the source regions 116a and 118a, and the drain regions 116b and 118b. \nReferring to FIG. 1C, an SAB layer 122 is formed over the substrate 100, covering the gate structures 106 and 108, and the exposed surface of the substrate 100.\nReferring to FIG. 1D, the SAB layer 122 in the peripheral circuit area 104 is removed, and the SAB layer 122a in the device area 102 is reserved. A metal layer 124 is then formed over the substrate 100, covering the SAB layer 122a, the gate electrode 108, and the exposed surface of the substrate 100.\nReferring to FIG. 1E, a thermal process is performed so that a portion of the metal layer 124 reacts with silicon under the metal layer 124 to form a metal silicide layer 126. The unreacted metal layer 124 is then removed.\nIn the process described above, the formation of the metal silicide layer can solve the problem of high resistance of the device. However, when the size of the device shrinks, the lattice distance in the channel region 128 seriously affects the electron mobility. The lattice distance becomes an essential factor in determining the device performance."}
{"text": "Online video chatting, video conferencing, video phone, wireless video communication, and/or so on, are very popular forms of communication. Advancements in video communication technology have enabled real-time interaction and sharing of information and documents over distances. This combination of video, voice and data creates a collaborative environment that can nearly replicate a face to face meeting. Adding digital effects to a video stream makes video communication more fun, functional, aesthetic, commercial, and/or serves other purposes."}
{"text": "1. Field\nExample embodiments generally relate to fuel structures and materials used in nuclear power plants.\n2. Description of Related Art\nGenerally, nuclear power plants include a reactor core having fuel arranged therein to produce power by nuclear fission. A common design in U.S. nuclear power plants is to arrange fuel in a plurality of fuel rods bound together as a fuel assembly, or fuel bundle, placed within the reactor core. These fuel rods typically include several elements joining the fuel rods to assembly components at various axial locations throughout the assembly."}
{"text": "The present invention relates to a knockdown display stand for use in merchandising products, and more particularly to a portable display stand designed for compact storage and easy relocation from one store to another.\nIn the merchandising of a number of products such as packaged, canned or bottled food or drink products, the products are typically displayed on the tiered shelves of a display stand at a supermarket or other retail store. Display stands or racks of such a type, frequently, do not belong to retail stores, but belong to the wholesalers or the manufacturers of the displayed merchandise. As a result, relocation of a display stand from one store to another is often required to efficiently utilize the same display stand in number of stores. Also, since the number of items displayed at a store fluctuates from week to week, storage of additional or unnecessary shelves at the store may be required. For convenience of relocation and storage, knockdown or collapsible structure of a display stand has been an important factor.\nAmong the prior art, U.S. Pat. Nos. 3,121,496 and 3,523,694 disclose display stands which can be knocked down or collapsed easily into compact units, respectively. These display stands include basket-type shelves, each being of a preassembled structure having built-in posts or spacer members adapted to be vertically extended between two shelves. Due to their outdated appearance and inefficiency in the manufacturing, display stands of these types are not in widespread use.\nU.S. Pat. Nos. 4,467,927 and 4,930,643 describe knockdown display stands formed, typically, of molded plastics. These display stands include tray-shaped shelves and separate posts designed to be positioned between the shelves. Each shelf is provided with integrally formed sockets for releasably receiving the ends of the posts, and the display stand may be conveniently taken apart. One disadvantage of such display stands is that upon disassembly, it is not easy to keep the parts for a display stand together in organized condition. This can result in loss of some of the parts, the posts in particular, of a display stand during \"knock down\" relocation, or \"knock down\" storage. Further, a box or other container must be found for holding the parts.\nWhat is needed, therefore, is a display stand for merchandising products. Such a display stand should be efficiently manufacturable, capable of being knocked down into a compact unit, and easy to keep in organized condition upon disassembly."}
{"text": "Mobile platforms are designed to service traffic from various bus mastering devices such as, for example, an USB-connected Bluetooth radio, with low response time, typically less than a few microseconds. As a result, most components of mobile platforms remain in a high state of readiness even when devices are idle or otherwise underutilized.\nEnergy efficiency is a key consideration to improve battery life for mobile platforms. Devices play a key role in energy efficiency for mobile platforms. The device performance depends on a central processor unit (CPU) as well as hardware and software system components such as, for example, the operating system, the graphics and I/O subsystems, application software, and/or host memory.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "Electrical connectors for use in connecting coaxial transmission lines to a variety of signal sources are known and have been widely used for decades. One particular application of such connectors involves connecting a signal receiving antenna to a coaxial transmission line, and one such known communication signal receiving arrangement 10 is shown in FIG. 1. Referring to FIG. 1, signal receiving arrangement 10 includes a tower 12 having an antenna 14 affixed thereto, wherein antenna 14 may be a dipole, microwave or other known antenna type. A coaxial electrical connector 16 is connected at one end to a matingly configured electrical connector coupled to antenna 14, and at its opposite end to a coaxial transmission line 18. Transmission line 18 is typically routed to a base station 20 located adjacent to tower 12 where signal evaluation and diagnostics are performed before providing the received signal to users via an output coaxial transmission line 22.\nConnector 16 and antenna 14 each define complementarily or matingly configured electrical connection structures such that connector 16 may be electrically connected to antenna 14 in a structurally secure manner. An example of one such connector 16 is described in U.S. Pat. No. 5,595,502 to Allison, which is assigned to the assignee of the present invention, and the disclosure of which is incorporated herein by reference. In this embodiment, the electrical connector coupled to antenna 14 is matingly configured to receive the '502 connector in both an electrically contacting and structurally secure manner.\nThe present invention provides a device for connecting and sealing between a coaxial transmission line connector and a signal source connector, wherein the device is configured to prevent liquid ingress.\nThe present invention further provides for an electrical connector having a liquid sensor disposed therein producing a signal indicative of the liquid content within the connector.\nThese and other objects of the present invention will become more apparent from the following description of the preferred embodiment."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus such as a xerographic type of copying machine or printer and, more specifically, to an image forming apparatus which has transporting means for transporting a recording sheet between a transfer section for transferring a toner image to the recording sheet and a fixing section for fixing the toner image transferred to the recording sheet and is capable of changing the recording-sheet transporting speed of the transporting means.\n2. Description of the Related Art\nA conventional type of image forming apparatus such as a xerographic type of copying machine or printer is constructed to transfer a toner image formed on a photoconductor drum to a recording sheet by means of a transfer section, transport the recording sheet to which the toner image was transferred to a fixing section by means of transporting means, and fuse and fix the toner image to the recording sheet by heat and pressure by means of a fixing unit provided in the fixing section, thereby forming an image.\nIn the fixing unit, since a toner image is fused and fixed to a recording sheet by heat and pressure, it is necessary to set its fixing speed according to a fixing characteristic due to the thickness of the recording sheet or the like or a fixing characteristic due to the kind of toner image. Specifically, in the case of a recording sheet or toner image which is difficult to fix, it is necessary to set the fixing speed to a slow speed, while in the case of a recording sheet or toner image which is easy to fix, it is necessary to set the fixing speed to a fast speed. In this respect, the fixing unit differs from the transfer section which electrostatically transfer a toner image formed on the photoconductor drum, at an approximately constant transfer speed. For this reason, in general, the time required for fixing in the fixing unit, i.e., the recording-sheet transporting speed (fixing speed) in the fixing unit, differs from the recording-sheet transporting speed (transfer speed) in the transfer section. Accordingly, when a recording sheet is to be transported from the transfer section to the fixing section, it becomes necessary to adjust the transporting speed of the recording sheet.\nTo adjust the transporting speed of the recording sheet, it is more desirable that the transporting path from the transfer section to the fixing section be made longer, because a speed difference between the transfer speed and the fixing speed is easier to be absorbed and a sufficient fixing time can be ensured.\nOn the other hand, to reduce the size of the image forming apparatus, it is more desirable to make the recording-sheet transporting path shorter.\nAn art for reducing the size of the apparatus while maintaining optimum fixing speed conditions is disclosed in, for example, Japanese Patent Laid-Open No. 171277/1997. This specification proposes a recording medium transporting system for a xerographic apparatus which defines the relationships between a distance l.sub.1 from a transfer position to the entrance-side axis of a transfer belt, a distance l.sub.2 from the entrance side to the exit side of the transfer belt, a distance l.sub.3 from the exit side of the transfer belt to a speed change point, a distance l.sub.4 from the speed change point to a fixing roller, a maximum length A.sub.max of a usable recording sheet, a minimum length A.sub.min of a usable recording sheet, a transfer speed V.sub.1 and a fixing speed V.sub.2, thereby making it possible to minimize the size of a transporting system which leads to a fixing section from the toner image transfer position, particularly, the length of the transporting belt or the distance to a nip position of the fixing roller.\nHowever, this proposal has the problem that the size, transfer speed, fixing speed and transfer position of a recording medium to be used as well as the length, fixing position and the like of a transporting unit impose restrictions on one another, thereby hindering free layout of individual units."}
{"text": "1. Field of the Invention\nThe present invention relates to a dryer, and more particularly, to an apparatus for opening/closing a door of a dryer.\n2. Discussion of the Related Art\nGenerally, a dryer is an apparatus for drying a washed laundry quickly using hot air. Lately, a washing machine equipped with a drying function has been developed as well.\nFIG. 1 is a perspective view of a dryer according to a related art.\nReferring to FIG. 1, a dryer according to a related art consists of a case 2, a drum 4, a heater (not shown in the drawing), and a motor (not shown in the drawing).\nAn opening 2a via which a laundry is put in or pulled out of the drum 4 is formed at a front side of the case 2. The drum 4 is rotatably installed in the case 2 to communicate with the opening 2a. A plurality of lifters 3 protrude from an inner conference of the drum 4 to lift the laundry in the drum 4 to fall.\nThe heater (not shown in the drawing) generating hot air for drying the laundry and the motor (not shown in the drawing) for rotating the drum 4 are installed outside the drum 4.\nThe case 2 consists of a cabinet 32, a front panel 33, a top cover 35, and a control panel 36.\nThe front panel 33 having the opening 2a formed at its front side is provided to a front side of the cabinet 32. The top cover 35 is provided to a top of the cabinet 32. The control panel 36 is provided on a rear top of the top cover 35. And, electronic units for controlling the dryer are installed in the control panel 36.\nMeanwhile, a door 10 is provided at one side of the opening 2a in the front panel 33 to open/close the opening 2a. \nThe door 10 consists of a hook 30 and a latch assembly 34. The hook 30 is fixed to one side of the door 10 using screws 31, and the latch assembly 34 is fixed to one side of the front panel 33 using screws 34a. And, the door 10 is provided to a hinge 12 fixed to the front panel 33 in the vicinity of the opening 2a to revolve centering around the hinge 12.\nThe latch assembly 34 locks or unlocks the hook 30 installed at the door 10 according to a signal of a control unit. In this case, a user freely enables to open/close the door 10 while power is not applied to a door lock switch separately installed at one side of the dryer.\nMeanwhile, while the door 10 is closed and the power is applied to the door lock switch, the latch assembly 34 keeps locking the hook 30 so that the door 10 is unable to be opened.\nHowever, the related art dryer has the following problem or disadvantage.\nFirst of all, the dryer differs from a washing machine in that water is not held in the dryer. Hence, the user needs to open the door in order to put more dry objects in the dryer in progress or in order to check a drying state of the laundry while the dryer operates. In this case, the structure of the apparatus for opening/closing the door of the dryer is equal or similar to that of the washing machine. Hence, it is inconvenient or takes pains to open/close the door."}
{"text": "The invention relates generally to transfer pad printing systems, and more particularly to pad printing plates and methods therefor.\nTransfer pad printing is a process of transferring ink images formed on etched printing plates, known as cliches, onto products with printing pads. Transfer pad printing is particularly suitable for printing on non-flat or irregular surfaces, for example on golf balls.\nOne known type of pad printing plate comprises a hardened and relatively thick steel plate, for example about 10 mm thick with a Rockwell C scale hardness of approximately 65. A surface of the steel plate is ground and polished, or lapped, to provide a relatively flat and smooth surface on which an image is formed or etched by means well known in the art.\nThe referenced hardened and relatively thick steel pad printing plates have the advantage of being very durable with an operating life expectancy of over one million or more printing cycles. These hardened thick steel pad printing plates however are costly to manufacture, particularly the grinding and lapping operations performed thereon. The plates also have a very high raw material cost. Moreover, these plates are heavy and bulky, and thus difficult to handle and costly to ship and store.\nAnother known type of pad printing plate comprises a relatively thin metal foil disposed on a base substrate. Thin metal foil plates however are relatively soft, with a Rockwell C scale hardness of about 44, and thus have a comparatively limited life expectancy.\nAn object of the present invention is to provide novel transfer pad printing systems, pad printing plates, combinations thereof and methods therefor that improve upon and overcome problems in the art.\nAnother object of the invention is to provide novel transfer pad printing systems, pad printing plates, combinations thereof and methods therefor that are economical.\nA further object of the invention is to provide novel transfer pad printing systems and pad printing plates therefor having reduced weight and relatively good durability.\nA more particular object of the invention is to provide novel transfer pad printing systems and pad printing plates therefor comprising a metal plate having a Rockwell C scale hardness not less than approximately 55 and a thickness between approximately 0.015 inches and approximately 0.020 inches.\nAnother more particular object of the invention is to provide novel transfer pad printing plates and methods therefor comprising etching an image on an upper surface of a metal plate having a Rockwell C scale hardness not less than approximately 55 and a thickness between approximately 0.015 inches and approximately 0.020 inches, and mounting a lower surface of the metal plate on a mounting surface of a base plate.\nThese and other objects, aspects, features and advantages of the present invention will become more fully apparent upon careful consideration of the following Detailed Description of the Invention and the accompanying Drawings, which may be disproportionate for ease of understanding, wherein like structure and steps are referenced generally by corresponding numerals and indicators."}
{"text": "The present invention relates to wireless communications, and in particular, to a unique technique for controlling amplifier performance in multiple communication modes by varying a supply voltage for the amplifier.\nThe present invention is particularly useful in mobile terminals, such as personal communication assistants, pagers, headsets, wireless modems, analog and digital cellular telephones, and the like. Since many of these devices are battery-powered, amplifier efficiency is preferably maximized to extend battery life. When amplifiers are designed for their highest efficiency in converting DC energy into RF energy, parasitic losses are minimized, bandwidths are reduced to a bare minimum, harmonics are terminated, and high-Q matching networks are employed. Unfortunately, these design goals are counter to current approaches used to implement multi-mode amplifiers capable of operating at multiple frequencies.\nMost RF power amplifiers are designed to operate over a single band of frequencies. If coverage is desired for one or more additional frequency bands, a multi-band amplifier is typically created. For a dual-mode application, the most straightforward approach is to simply use two amplifiers and switch between them to select a desired band for transmission. If a single amplifier configuration is desired, then matching networks that provide the proper impedance transformation for both frequency bands are required. These matching networks are implemented in a number of ways. The use of series and parallel resonant elements, typically inductors and capacitors, in the matching networks are selected such that at one frequency band, the combination appears inductive and at the other band, it appears capacitive. Clever combinations of series and shunt element pairs may allow a creation of networks that deliver the desired properties over multiple frequency bands. These approaches typically result in limited bandwidth for each band and difficulty in tuning and maintaining performance over each of the bands.\nIf the amplifier is designed to operate in one band at a time, one or more switches are typically used to add elements to or remove elements from the matching networks for the appropriate band. This works well and has been employed in many applications. The drawbacks are the additional energy needed to operate the switches and the losses the switches add to the networks. If a change in the amplifier mode is desired, such as changing linearity, efficiency, or power, then the load of the amplifier may be switched by adding or removing elements from the matching network. Examples of these techniques are disclosed in U.S. Pat. Nos. 5,438,684 and 5,673,287, which are assigned to Motorola, Inc. Again, the switching techniques reduce system efficiency, which results in decreased battery life. As such, there is a need for an improved and efficient multi-mode amplification technique that does not require inefficient and complicated matching networks and amplifier designs.\nThe present invention eliminates the need for complex impedance networks or parallel amplification stages for multi-mode mobile terminals. A wideband power amplifier is configured to amplify signals in different frequency bands corresponding to different operating modes. The supply voltage of the wideband power amplifier is adjusted in light of the load impedance of radiating circuitry to achieve a desired output power for the respective operating modes.\nAccordingly, the present invention relates to a radio frequency communication system including radiating circuitry, a wideband power amplifier output stage, a variable power supply, and a control system. The radiating circuitry, which includes an antenna, has a load impedance. The wideband power amplifier output stage is coupled to the radiating circuitry and is adapted to amplify a radio frequency signal in each of a plurality of frequency bands corresponding to a plurality of operating modes. The variable power supply is adapted to provide a selected supply voltage for supplying power to the wideband power amplifier output stage in response to a power control signal. The control system is adapted to generate the power control signal based on a selected one of the plurality of operating modes. The selected supply voltage and the load impedance for the frequency band of the selected operating mode determine the power delivered to the radiating circuitry from the wideband power amplifier output stage for each of the plurality of operating modes.\nThe power level for any given operating mode may be further controlled by controlling a signal level for the radio frequency signals to be amplified by the wideband power amplifier. Alternatively, controlling the selected supply voltage may control the power level.\nThe wideband power amplifier output stage is preferably made of transistors configured to operate in a saturation region when in one of the selected operating modes and in a linear region when in another of the selected operating modes. Preferably, one or more wideband intermediate amplifier stages are coupled in series to the wideband power amplifier output stage. The intermediate amplifier states are fed by combining circuitry for passing signals within the frequency band of the selected operating mode.\nBias circuitry is used to provide bias to the wideband power amplifier output stage. The bias is configured to optimize efficiency of the wideband power amplifier output stage at different supply voltages for each of the operating modes. Further, the radiating circuitry may include a selectable impedance component to adjust the reactance of the load impedance for at least one of the operation modes to optimize signal transmission. The radiating circuitry may be configured to provide a load impedance having a first impedance for a first frequency band and a second impedance for a second frequency band. The first and second impedances may be substantially the same or may vary based on design considerations. In operation, the supply voltages for the wideband power amplifier output stage will set the output power levels.\nThose skilled in the art will appreciate the scope of the present invention and realize additional aspects thereof after reading the following detailed description of the preferred embodiments in association with the accompanying drawing figures."}
{"text": "Security fences of various types are known in the patent literature and available on the market. These include taut-wire fence systems which employ tensionsed wires mounted onto motion sensors, whereby an attempt to cut or spread apart the wires results in an alarm indication of the approximate location of the attempted intrusion.\nThere are also known a variety of security barriers which employ a fiber optics sensing apparatus. U.K. published patent application No. 2,098,770 describes a security barrier structure comprising a lattice of hollow tubular members through which fiber optic cable is threaded. An attempt to break through the barrier breaks or distorts the fiber by overtensioning same, thus causing a sensible attenuation of the an optical signal transmitted through the cable.\nU.K. published patent application Nos. 2,038,060; 2,046,971; and 2,062,321 and U.S. Pat. Nos. 4,292,628 and 4,399,430 all show security applications, wherein an alarm indication is provided by breakage of an optical fiber. U.K. published patent application No. 2,077,471 shows a security application wherein a pressure sensitive fiber optic composite cable is provided."}
{"text": "Venomous bites and stings can cause a variety of reactions depending upon the source of the venom and the sensitivity of the individual or animal. In some cases, venom from a bite or sting can cause anaphylaxis, an immediate hypersensitivity which can be life-threatening. In other cases, certain venoms can cause cutaneous \"local\" reactions. Cutaneous local reactions can be characterized as 1) \"non-allergic\" reactions which are of limited size and duration or as 2) \"allergic\" or \"large\" local reactions which are typically larger in size and longer in duration. With regard to Hymenoptera venoms (bee, wasp, hornet and yellowjacket), \"the non-allergic local reaction is a toxic response to venom constituents, while the large local reaction appears to be caused by an allergic reaction to venom proteins.\" See, D. N. Wright, Local Reactions To Stinging Insects (Hymenoptera), Allergy Proc. 11(1):23-28 Jan.-Feb. 1990).\nUpon receiving a venomous bite or sting, a variety of symptoms can be exhibited due to the venom, including pruritus, erythema, urticaria, angioedema, soft tissue swelling, inflammation of the affected area and pain in the affected area.\nWhen injected subcutaneously, many venoms from bites and stings induce extravasation from adjacent blood vessels. See, V. Cattell, Focal Mesangial Proliferative Glomerulonephritis In The Rat Caused By Habu Snake Venom: The Role Of Platelets, British J. of Exp. Pathol., 60(2):201-208 (April 1979); In addition, the venoms induce platelet aggregation and mast cell degranulation, two components of inflammation along with extravasation. Serotonin release has also been associated with the injection of venoms. See, for instance, Y. Ozaki, et al, Mastoparan, A Wasp Venom, Activates Platelets Via Pertussis Toxin-Sensitive GTP-Binding Proteins, Biochem. Biophys. Res. Commun., 170(2):779-85 (Jul. 31, 1990) and C. Wang, et al, Experimental Study of Chinese Agkistrodon Actus Venom In Activation Of Rabbit Platelets In Vivo, Hua Hsi I Ko Ta Hsueh Hsueh Pao, 25(1):38-40 (March 1994).\nSince the discovery of serotonin (5-hydroxytryptamine, 5-HT) over four decades ago, the cumulative results of many diverse studies have indicated that serotonin plays a significant role in the functioning of the mammalian body, both in the central nervous system and in peripheral systems as well. Morphological studies of the central nervous system have shown that serotonergic neurons, which originate in the brain stem, form a very diffuse system that projects to most areas of the brain and spinal cord. R. A. O'Brien, Serotonin in Mental Abnormalities, 1:41 (1978); H. W. M. Steinbusch, Handbook of Chemical Neuroanatomy, Volume 3, Part II, 68 (1984); N. E. Anden, et al., Acta Physiological Scandinavia, 67:313 (1966). These studies have been complemented by biochemical evidence that indicates large concentrations of 5-HT exist in the brain and spinal cord. H. W. M. Steinbusch, supra.\nWith such a diffuse system, it is not surprising that 5-HT has been implicated as being involved in the expression of a number of behaviors, physiological responses, and diseases which originate in the central nervous system. These include such diverse areas as sleeping, eating, perceiving pain, controlling body temperature, controlling blood pressure, depression, schizophrenia, and other bodily states. R. W. Fuller, Biology of Serotonergic Transmission, 221 (1982); D. J. Boullin, Serotonin in Mental Abnormalities 1:316 (1978); J. Barchas, et al., Serotonin and Behavior, (1973).\nSerotonin plays an important role in peripheral systems as well. For example, approximately 90% of the body's serotonin is synthesized in the gastrointestinal system, and serotonin has been found to mediate a variety of contractile, secretory, and electrophysiologic effects in this system. Serotonin may be taken up by the platelets and, upon platelet aggregation, be released such that the cardiovascular system provides another example of a peripheral network that can release and respond to serotonin. Given the broad distribution of serotonin within the body, it is understandable that tremendous interest in drugs that affect serotonergic systems exists. In particular, receptor-specific agonists and antagonists are of interest for the treatment of a wide range of disorders, including anxiety, depression, hypertension, migraine, compulsive disorders, schizophhrenia, autism, neurodegenerative disorders, such as Alzheimer's disease, Parkinsonism, and Huntington's chorea, and cancer chemotherapy-induced vomiting. M. D. Gershon, et al., The Peripheral Actions of 5-HYDROXYTRYPTAMINE, 246 (1989); P. R. Saxena, et al., Journal of Cardiovascular Pharmacology, 15:Supplement 7 (1990).\nSerotonin produces its effects on cellular physiology by binding to specialized receptors on the cell surface. Multiple types of receptors exist for many neurotransmitters and hormones, including serotonin. The existence of multiple, structurally distinct serotonin receptors has provided the possibility that subtype-selective pharmacologic agents can be produced. The development of such compounds could result in new and increasingly selective therapeutic agents with fewer side effects, since activation of individual receptor subtypes may function to affect specific actions of the different parts of the central and/or peripheral serotonergic systems.\nAn example of such specificity can be demonstrated by using the vascular system as an example. In certain blood vessels, stimulation of certain 5-HT receptors on the endothelial cells produces vasodilation while stimulation of certain 5-HT receptors on the smooth muscle cells produces vasoconstriction.\nCurrently, the major classes of serotonin receptors (5-HT.sub.1, 5-HT.sub.2, 5-HT.sub.3, 5-HT.sub.4, 5-HT.sub.5, 5-HT.sub.6, and 5-HT.sub.7) contain some fourteen to eighteen separate receptors that have been formally classified based on their pharmacological or structural differences. For an excellent review of the pharmacological effects and clinical implications of the various 5-HT receptor types, see Glennon, et al., Neuroscience and Behavioral Reviews, 14:35 (1990).! discoveries.\nOne class of serotonin receptors is the 5-HT.sub.2. Of this class, several subtypes are known to exist. These subtypes include 5-HT.sub.2A, 5-HT.sub.2B and 5-HT.sub.2C. The subtype 5-HT.sub.2A is located in many tissues, including but not limited to, the vascular smooth muscle, platelets, lung, CNS and gastrointestinal tract. This receptor is thought to be associated with several effects: for example, vasoconstriction, platelet aggregation, and bronchoconstriction. The 5-HT.sub.2B receptor is localized in the rat lung, stomach fundus, uterus, bladder, and colon. Interesting areas of 5-HT.sub.2B receptor localization in the human include, but are not limited, to the brain and blood vessels. Subtype 5-HT.sub.2C is located in the CNS with a high density in the choroid plexus.\nBecause of the widespread dissatisfaction with the current treatments for venomous bites and stings within the affected population, there exists a need for a more efficacious and safe treatment. The present invention provides such a treatment."}
{"text": "Stair lifts are well known and are particularly used for transporting loads such as persons and/or goods up or down over a stairs by being carried on a chair or a platform. In general a stair lift comprises of a guide or rail to be attached along a staircase and a motorized carriage or trolley which carries and guides the load along the rail.\nStair lifts may be driven by friction drive rollers which are in frictional engagement with the guide. These have proven to be less expensive than form fitted drive means such as sprocket and chains or rack and pinion.\nFor example, a stair-climbing device is known from patent publication EP-0,881,188-A1. This publication shows a stair-climbing device comprising a guiding rail and a trolley movable on the guiding rail. The stair-climbing device is further provided with a motor that drives a drive roller supported by the trolley, wherein the drive roller is in frictional engagement with the guiding rail. Each longitudinal side of the guiding rail is in frictional engagement with a guiding roller to keep the trolley in a predetermined position with respect to the guided rail. Each drive roller is provided with a means for pressing the drive roller into frictional engagement.\nDrawback of these known stair lifts is that they are mechanically complex."}
{"text": "Glycopyrrolate, the active pharmaceutical ingredient in Robinul® tablets, Robinul® Forte tablets and Robinul® injection, is a quaternary ammonium compound having the chemical name 3-[(cyclopentylhydroxyphenylacetyl)oxy]-1,1-dimethylpyrrolidinium bromide. Glycopyrrolate is an anticholinergic and antispasmodic that inhibits gastrointestinal nerve receptor sites that stimulate both the secretion of stomach acid and smooth muscle activity in the digestive tract. Accordingly, it diminishes the volume and free acidity of gastric secretions and controls excessive pharyngeal, tracheal, and bronchial secretions.\nGlycopyrrolate is indicated for use as adjunctive therapy in the treatment of peptic ulcer. See PHYSICIANS' DESK REFERENCE (57th ed., Medical Economics Co., 2003). The most frequent adverse reactions to glycopyrrolate include those that are common to anticholinergics generally including, for example, xerostomia, decreased sweating, urinary hesitancy and retention, blurred vision, tachycardia, palpitations, headaches, dizziness, nausea, vomiting, nervousness, and other reactions.\nPharmacokinetic studies have not previously been conducted to evaluate the effect of food on the pharmacokinetics of glycopyrrolate. It is desirable to increase the oral bioavailability of a drug substance, such as glycopyrrolate, to increase the extent of the therapeutic effect on the user. In general, food has a variable effect on the bioavailability of an active agent. Drug-food interactions may result in reduced, delayed or increased systemic drug availability. See, e.g., Clin Pharmacokinet 1984 September-October; 9(5):404-34. Food may interact with a coadministered drug at the following phases: (i) before and during gastrointestinal absorption; (ii) during distribution; (iii) during metabolism; and (iv) during elimination. See Clin Pharmacokinet 1999 September; 37(2):213-55. For some drugs such as, for example, nefazadone and troglitazone, bioavailability increases with food. See Clin Pharmacokinet 1997 October; 33(4):260-75; Clin Pharmacokinet 1999 August; 37(2):91-104.\nThere remains a need for a method of treatment of peptic ulcer and other gastrointestinal disorders using glycopyrrolate that has increased oral bioavailability and, thus, improved efficacy following administration to human patients."}
{"text": "The present invention relates to goal apparatus, in particular a portable goal apparatus comprising a frame and first and second flexible goal post members and a flexible crossbar member that are releasably connectable to the frame.\nMany sports such as football (soccer), hockey etc. utilize a structure having parallel posts and a crossbar to define a goal area. Such structures are usually permanent or semi permanent fixtures and are made from steel or other heavy and cumbersome materials. However, it is also known to provide portable goals that may be assembled on site prior to use and disassembled once use has ceased. In this manner, a location need not be permanently devoted to use as a sports field. Portable goals are particularly useful in training scenarios where it may be beneficial to have a large number of goals to train with, which can then be dismantled after use.\nPortable goals known in the prior art often involve the use of tubular plastic materials to form the posts and cross bar which, while representing a saving on weight, are still very bulky when disassembled and therefore represent a significant problem with regard to carrying and storage.\nIt is an aim of aspects of the present invention to address the above mentioned or other problems."}
{"text": "Shrews are a very ancient group of primitive mammals that resemble most closely the proto-mammals. They are not closely related to rodents because rodents evolved from different groups of mammals. According to Dufton (1992), the known venomous species of shrew are: the northern short-tailed shrew (Blarina brevicauda), the Haitian solenodon (Solenodon paradoxus), the European water shrew (Neomys fodiens) and the Mediterranean shrew (Neomys anomalous). Another venomous shrew is the southern short-tailed shrew (Blarina carolinensis). It has also been suggested that the Cuban solenodon (Apotogale cubanus), the American shrew (Sorex cinereus) and the Maritime shrew (Sorex maritimensis) could be venomous. The northern short-tailed shrew (Blarina brevicauda) and its closely related species use a paralytic venom in its saliva to paralyze insects, other invertebrates (worms, annelids etc.), nesting birds and small mammals which it then stores, alive in its den, for future feeding (Martin 1981; George et al. 1986; Dufton 1992).\nThe shrew venom literature generally consists of seven articles from the 40s and 50s and one MA thesis in 1966 [Christenbury 1966]. These are summarized in a review [Dufton 1992]. Using a crude ammonium sulfate precipitate of shrew saliva glands, Ellis and Krayer (1955) concluded the active agent was probably a protein and, because of its inability to dialyze, a larger protein. A major contribution of the Ellis & Krayer work was to show activity in cats, dogs, mice, rats, guinea pigs and rabbits. Christenbury [1966] showed Ellis & Krayer's preparation stopped oxygen consumption by mouse kidney and liver slices. Japanese patent application (JP 10-236963; 1998) appears to disclose an alcoholic extract of saliva glands from two shrew species (Sorex unguiculatus & Sorex shinto saevus) as a calcium channel blocker and its use as a hypotensive. The purity is low—the extract includes any compounds that would dissolve in 70% ethanol. There is no information about the responsible active molecule/s in the unknown mixture of compounds."}
{"text": "Because of the extreme nature of the application, materials used for fabrication of ammunition cartridges must demonstrate excellent mechanical and thermal properties. The prevalent materials for production of cartridge casings for all calibers of ammunition in the world today are metals. Brass is the leading material, followed in smaller amounts by steel and, in limited amounts, aluminum. The use of polymeric materials for ammunition cartridge casings has been extensively investigated over the past 40 years, but success has been elusive.\nBrass, steel, and, to a lesser degree, aluminum cartridge casings suffer from a number of disadvantages, the most important of which are heavy weight and corrosion concerns. Aluminum has an added disadvantage of potential explosive oxidative degradation and is thus used only in low-pressure cartridges or in applications that can tolerate relatively thick casing walls. Given these issues, desirable materials for ammunition cartridge casing fabrication would be lightweight and impervious to corrosion while having mechanical properties suitable for use in ammunition applications. Many lightweight polymeric materials are sufficiently corrosion resistant; however, to date, polymers have been used only in niche ammunition applications where their inferior mechanical and thermal properties can be tolerated (e.g., shotgun shells contain polyethylene components).\nWhile stability under broad ranges of handling and storage conditions are crucial, the greatest mechanical demands on the cartridge material are experienced during the firing event. The material at the cartridge base end, which supports the primer, must first absorb the impact of a firing pin on the primer without mechanical failure. Upon ignition and combustion of an encapsulated propellant, rapidly expanding gases create high pressure, which expels a projectile from the barrel of the fired firearm. The ammunition cartridge casing must withstand and contain the pressure developed by the explosion so that the gaseous combustion products expand only in the direction of the barrel opening, thus maximizing energy conversion to projectile kinetic energy.\nA firearm's cartridge chamber closely fits the outside of a cartridge and thus supports the majority of the cartridge casing wall in the radial direction; however, in many firearms, a portion of the cartridge base end protrudes from the chamber and is thus unsupported. During firing, a stress profile is developed along the cartridge casing, with the greatest stresses being concentrated at the base end. Therefore, the cartridge base end must possess the greatest mechanical strength, while a gradual decrease in material strength is acceptable axially along the casing toward the forward end which receives the projectile.\nA typical brass cartridge casing is engineered to provide a strength profile along the casing length which reflects the varying mechanical demands, with the strongest and hardest material located at the cartridge base end. In brass and other metals, a strength profile is easily induced by varying the heat treatment conditions from one end of the casing to the other, but this is not an option for polymers. A mechanical strength profile can be achieved in a polymeric ammunition cartridge casing by varying the casing wall thickness; however, where the casing external geometry is fixed by existing firearm chamber size, an increased casing wall thickness often results in a casing with insufficient internal volume to accept the required propellant charge.\nMany ammunition articles have been designed with cartridge cases comprised of two or more separate parts. The individual components are typically fabricated from different materials; a high-strength, usually metallic material comprises the cartridge casing base portion or “cap” while a polymer or other material comprises the remainder of the casing. For example, commercial shotgun ammunition employs a metallic base or cap joined to a polymeric top or sleeve. Although a significant amount of metal is required for such an ammunition cartridge, weight and cost savings can be sufficient to make it commercially acceptable.\nWhile the most severe mechanical requirements of an ammunition cartridge are focused on the base end, the top or forward portion of the casing must meet several material requirements as well. Upon combustion of the cartridge propellant, a very large quantity of energy is released in a matter of a few milliseconds, thus producing very high stresses and strain rates. The casing material must possess adequate ductility to absorb the shock of the explosion without experiencing brittle fracture. Also, the material must possess sufficient rigidity and strength to avoid creep, flow, or other deformation.\nA vast amount of effort has been dedicated to designing plastic ammunition cartridges, and researchers in the field have tested a variety of materials. Despite these attempts, consistent success has not been achieved.\nBecause of the demands on the casing material, the key problem in developing polymer-cased ammunition remains identifying a suitable polymeric material. It appears that all of the polymeric materials tried thus far are critically deficient in either their absorption of the impact energies generated during the firing event or their retention of mechanical integrity at high temperatures. A significant improvement in the art would be the identification of polymeric materials capable of at least serving as the top or forward portion of the ammunition cartridge casing."}
{"text": "1. Field of the Invention\nThe invention relates to water bed mattress liners; and, more particularly, to a tackless liner for the mattress of a water bed which is installed without the need for securing the same to the bed frame.\n2. Description of the Prior Art\nIt is well known that a water bed mattress holds a considerable amount of water and has been known to leak or rupture in the past. In order to contain water that might leak or escape from the mattress, most water beds have a liner which is of generally water tight material, such as vinyl, and is installed in the frame of the bed prior to placement of the mattress therein. The mattress thus sits on and inside of the liner which will hold any water leaking or escaping from the mattress.\nMany variations of liners have been suggested in the past. In one prior art arrangement, a vinyl flexible liner in the form of a rectangular sheet having upright sides is placed inside of the bed frame and the upright sides are nailed to the inside of the upright sides of the bed frame by nails passing through rigid strips of material thus securing the liner to the frame. Needless to say, such nails or any similar fastening means can work loose and make holds in the mattress causing damage thereto.\nIn another prior art arrangement, elongated strips of a rigid material are secured along the inside walls of the upright sides of the bed frame so that their upper surfaces are unrestricted. A flexible liner is then laid inside of the frame and the ends thereof are inserted over the unrestricted sides of the strips thereby trapping the liner between the strips and the frame. Such an arrangement requires hardware securing the strips to the frame, such as nails or screws, which can work loose and puncture the mattress. Further, the strips must be secured to the frame in a way such that they resiliently grasp or trap the liner so it won't shift or move out of engagement with the strips during use.\nBoth prior art arrangements discussed above obviously require elaborate installation and cannot easily be taken down if one desires to drain and move the water bed.\nThere is thus a need for a water bed liner which is easy to install, can be taken down quickly and easily, requires no nails, screws, tacks or hammering or other securing means and provides a perfect fit for the mattress."}
{"text": "In the related art, an image reading device determines a size of an original document using an image sensor used for reading the original document. The image sensor detects light emitted from a light source and reflected by the original document. Based on the intensity of the light detected by the image sensor, the image reading device determines the size of the original document.\nHowever, when a platen cover is not closed, the plurality of sensors detects ambient light coming from outside of the image reading device. Hence, in some cases, the image reading device may not precisely determine the size of the original document."}
{"text": "Amyotrophic lateral sclerosis (ALS), also known as Lou Gehrig's disease, is a debilitating disease with varied etiology. ALS is characterized by rapidly progressive weakness, muscle atrophy and fasciculations, muscle spasticity, difficulty speaking (i.e., dysarthria), difficulty swallowing (i.e., dysphagia) and difficulty breathing (i.e., dyspnea). Worldwide, ALS kills about 100,000 yearly, with 125,000 new ALS diagnoses yearly.\nALS causes muscle weakness and atrophy throughout the body due to the degeneration of the upper and lower motor neurons. Individuals affected by the disorder may ultimately lose the ability to initiate and control all voluntary movement, although bladder and bowel sphincters and the muscles responsible for eye movement are usually spared until the terminal stages of the disease. Most patients do not lose cognitive function, but about 5% develop frontotemporal dementia.\nThe precise cause of ALS is still unknown, but mutations in the gene that produces the Cu/Zn superoxide dismutase (SOD1) enzyme are present in approximately 20% of familial ALS cases. SOD1 is a powerful antioxidant that protects the body from damage caused by the toxic free radical superoxide, which is generated in the mitochondria. To date, over 110 different mutations in SOD1 have been linked with ALS. Failure of defenses against oxidative stress upregulates programmed cell death (apoptosis). An accumulation of free radicals due to faulty SOD1 functioning may thus lead to motor neuron degeneration. However, it is also possible that mutant SOD1 induces toxicity by a gain of function. For example, in some mutant SOD1 mice, mutant SOD1 aggregates (misfolded protein deposits) were found only in diseased tissues, and greater amounts were detected during motor neuron degeneration. This observation is consistent with a model where mutant SOD1 aggregates help disrupt cellular functions by damaging mitochondria, proteasomes or protein folding chaperones. In humans, SOD1 mutations cause only 2% of total ALS cases, and the etiological mechanisms may be distinct from those responsible for the sporadic form of the disease. For non-familial cases (around 90% of ALS cases), the disease may be triggered by head trauma, military service, or participation in contact sports.\nRiluzole (RILUTEK®) is the only drug currently approved to treat ALS, but its therapeutic success is limited. Riluzole slows down disease progression by about 3-9 months, prolongs survival by several months, and delays the need for ventilation support. However, the drug does not reverse the damage already done to motor neurons, and may cause liver damage in about 10% of the patients.\nIn one aspect, development of effective ALS treatments may be hampered by a combination of low drug bioavailability and disease-driven pharmacoresistance, which limits CNS drug penetration and ultimately compromises drug efficacy.\nPharmacoresistance is mediated for example by the ATP-binding cassette (“ABC”) drug efflux transporters. ABC transporters are transmembrane proteins that utilize the energy of adenosine triphosphate (ATP) hydrolysis to carry out biological processes, including translocation of substrates (e.g., metabolites, lipids, sterols and drugs) across extra- and intracellular membranes. ABC transporters play a crucial role in the development of multidrug resistance (MDR), whereby patients eventually develop resistance not only to the drug they are taking but also to additional types of drugs to which they had not been exposed. MDR is linked by the increased cellular excretion of the drug by ABC transporters. ABC transporter families include ABCB1 (also known as P-glycoprotein, P-gp or MDR1), which transports organic cationic compounds or neutral compounds; MRP, which transports organic anionic compounds; and ABCG2 (also known as breast cancer resistance protein or BCRP), which confers resistance to Topoisomerase I or II inhibitors such as topotecan, irinotecan, and doxorubicin. ABC transporters are predominantly located in the lumen of endothelial cells of the blood-brain barrier, where they limit the entry of neurotoxins into the CNS. Unfortunately, clinical trials in cancer patients to which ABC transporter inhibitors were co-administered, with the objective of increasing the anticancer drug efficacy, had ultimately disappointing outcomes (Sharom, 2008, Pharmacogenomics 9(1):105-127; Szakacs et al., 2006, Nat. Rev. Drug Discov. 5:219-234).\nThere is a need in the art for novel methods of treating neurodegenerative diseases in a mammal in need thereof. The present invention fulfills this need."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the hybrid variety designated CH977760, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include greater yield, better stalks, better roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen comes to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop hybrid parent plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants and have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
{"text": "Hazardous volatile organic compounds (VOC) in soils and groundwater can pose significant health risks, particularly if aquifers which feed the water supplies of population centers are threatened. Current remediation methods for treatment of soil and groundwater include air stripping, vacuum extracting, carbon containment, incineration/oxidation and bioremediation.\nAir stripping and vacuum extracting are dispersion technologies which extract vapors (the words vapors, VOCs and gases are generally used interchangeably throughout and, unless clearly understood otherwise, are intended to have the same meaning) from the soil and vent the contaminants to the atmosphere. While these techniques are relatively inexpensive and can prevent or reduce groundwater contamination threats, they are environmentally undesirable because they merely change the medium of the pollution (from soil to air). Venting is controversial and its use in the future is likely to be limited by stricter air quality regulations.\nCarbon containment is a collection technology in which vapors from vapor extraction wells or other VOC sources are passed through carbon (or other adsorbent) filled canisters to adsorb the VOC's. Collection technologies prevent or limit groundwater contamination by reducing the contaminant level and can be cost effective for certain flowrates and contamination levels. However, containment is not a destruction technology; the VOCs collected require removal and subsequent disposal. Also, containment is not universally effective, as some VOCs have low adsorptivity.\nIncineration involves high temperature burning of waste streams from soil or air-stripped ground water for destruction of VOCs. Incineration is highly controversial, often achieving only incomplete destruction, and is costly. Incomplete destruction can produce products more hazardous than the original contaminant and has the potential of releasing them to the atmosphere.\nThis current state-of-the-art has led to the present invention, which achieves safe and effective destruction on site for a wide variety of VOCs."}
{"text": "FIG. 8 shows a state in which sheets S are stacked in a sheet cassette serving as a sheet holding section of an image forming apparatus. FIG. 8 shows the interior of the sheet cassette from a side where the sheets S are fed. First, in FIG. 8A, an intermediate plate 401 on which the sheets S are stacked is provided in the sheet cassette, and the sheets S stacked on the intermediate plate 401 are regulated and held by side regulation plates 204 and 205. The side regulation plate 205 is movable in the direction of an illustrated arrow A, and regulates and aligns the sheets S in accordance with the size of the sheets S in a direction orthogonal to a conveying direction of the sheets S (also referred to as a widthwise direction of the sheets S). In addition, a leaf spring 206 for applying a predetermined pressure to the sheets S is provided at the side regulation plate 205 to regulate the sheets S. Further, at a side opposite to a side towards which the sheets are conveyed, the sheets S are also regulated and held by a rear-edge regulation plate, which is movable in accordance with the size of the sheets S.\nThe stacked sheets S are picked up by a pickup roller 103 and are fed into the image forming apparatus. If a predetermined number of sheets S is fed, and a sheet S is no longer detected by a sheet-surface sensor 104 that detects a topmost surface of the sheets S, the intermediate plate 401 is lifted up in the direction of arrow B until a sheet is detected by the sheet-surface sensor. Accordingly, while repeatedly picking up the sheets S and lifting up the intermediate plate 401, the sheets S are supplied. According to Patent Document 1, a stepping motor for lifting up stacked sheets S is used to repeat a lifting-up operation in accordance with the number of picked up sheets.\nFurther, in such a sheet cassette, the sheets S may not be properly fed even if the sheets S are picked up by the pickup roller 103. In such a case, the picking-up operation by the pickup roller 103 is repeatedly executed again to also perform control for reducing the occurrence of improper sheet-feeding (also called “sheet-feed retry control”). Patent Document 2 discusses a technology that reduces the occurrence of unsuccessful sheet-feeding by executing a picking-up operation again when a sheet S is unsuccessfully fed by a pickup roller.\nIn such a sheet cassette, when the environment in which the sheet cassette is used is a high-temperature and high-humidity environment, the sheets S absorb moisture, thereby reducing the toughness of the sheets. The sheets may become flexed or wavy by pressure of the side regulation plate 205. An example in which the sheets are wavy is shown in FIG. 8B. Accordingly, if the sheets absorb moisture and become wavy, a height of the sheets S opposing the pickup roller 103 (height indicated by an auxiliary line 1) is less than a height of the sheets at a position where the sheet-surface sensor 104 performs detection (height indicated by an auxiliary line 2), thereby producing a height difference (h). As mentioned above, on the basis of a detection result of the sheet-surface sensor 104, the intermediate plate 401 is lifted up to control the height of the sheets S. Therefore, when the sheets are wavy, pressure generated by contacting the pickup roller 103 with the sheets S is reduced, thereby causing an improper pickup operation.\nTo reduce the waviness of the sheets S, for example, the pressure of the side regulation plate 205 may be reduced. However, if the pressure of the side regulation plate 205 is reduced, the degree of alignment of the sheets S in the cassette is reduced. If the degree of alignment is reduced, the sheets S may be obliquely fed. If the sheets S are obliquely fed, improper sheet-feeding results. Therefore, this method of reducing the waviness of the sheets does not satisfactorily solve the problem.\nThe present invention is achieved for solving the aforementioned problems, and its object is to make it possible to reduce the occurrence of improper sheet-feeding when sheets are fed."}
{"text": "Such a tablet compressing machine is known from the document DE 20 2007 003 176 U1. The document DE 20 2007 003 176 U1 describes a tablet compressing machine, in which a cleaning liquid can be distributed, submerging the rotor and wherein ultra-sonic waves can be applied in the sealed housing to cleanse particles off from surfaces. Another tablet compressing machine is known from the document EP 0 637 507 A1. The document EP 0 637 507 A1 describes a method of washing a tablet compressing machine, wherein in a first step a cleaning liquid is distributed. In a second step, a flushing liquid is distributed and the liquid sucked out of the housing. In a third step, all openings of the housing are closed and a vacuum is applied to dry the sealed housing quickly.\nGeneric tablet compressing machines normally have a rotor unit comprising a die plate, an upper stamp holder and a lower stamp holder. The die plate is provided with a plurality of bores for inserting the upper or lower stamp during rotation. Thereby the upper and the lower stamps compress powdery material into shaped tablets. The upper and lower stamps are movably arranged in axial direction in associated upper and lower stamp holders. After the tablet has been compressed within a bore in the die plate, it is normally pushed out of the corresponding bore of the die plate by the lower stamps. The combination of upper stamps and holders, die plate and lower stamps and holders is called a rotor unit. Generic tablet compressing machines have a capacity of 500 to 25,000 tablets per minute depending upon the number of upper and lower stamps and the speed with which the rotor rotates.\nWhen changing to another batch, that is when the powdery material is changed, for introducing a new powder which is to be compressed, the previous powder must first be completely removed from the machine. No residues of the previous powder should get into the surrounding air, which is why the rotor units in tablet compressing machines are normally encapsulated in a sealed housing. Generic tablet compressing machines are usually equipped with a spraying apparatus by means of which a cleaning liquid such as water may be distributed for removing the residual medical substance. Powder components suspended in air are captured by the spray mist and bound by the liquid, whereupon they sink down and settle on horizontal surfaces of the rotor unit and its vicinity. Powder residues already laying on surfaces are bound in the same way. The cleaning liquid together with the bound powder residues forms puddles on horizontal surfaces which do not drain away. These surfaces may, for instance, be formed on the rotor unit, on ancillary units and parts of the housing. Usually the rotor unit and the spraying facility are at least partially disposed within the sealed housing so that no powder components can pass into the surrounding area during tablet production.\nIn order to drain the cleaning liquid off the tablet compressing machine it is known in the prior art to provide only inclined, that is non-horizontal surfaces within the tablet compressing machine, thus avoiding the forming of puddles. However, a lot of effort is required to provide only non-horizontal surfaces in such a complicated machine as a tablet compressing machine, in particular when a plurality of components are involved.\nThis is the state of the prior art on which the invention is based. The invention is aimed at reducing the work and effort involved in constructing and manufacturing a tablet compressing machine."}
{"text": "The proliferation of the Internet, in particular the World Wide Web (hereinafter the Web), makes a large amount of information accessible to anyone connected to the Internet. In order to access the Internet, a user typically subscribes to an on-line service provider. Using, for example, a modem coupled to a computer, the user contacts the service provider who, in turn, connects the user to the Internet. Once the user is coupled to the Internet, the user is said to be \"on-line\" (a user who is not on-line is referred to as being \"off-line\"). When on-line, the user is able to access numerous resources available through conduits such as, for example, the Web.\nThe Web consists, generally, of voluminous \"sites\". Various Web sites contain information of virtually innumerable type, style, and content. Web sites are often comprised of a plurality of Web documents referred to as Web pages. In general, the layout language for a Web document is Hypertext Markup Language (HTML). Web documents are stored in HTML format in association with \"server\" software on remote servers. A Web document is given a \"Uniform Resource Locator\" (URL) which is essentially an address path identifying the server which hosts the desired document plus the location of the document on the server. Using \"browser\" software, an end-user can send a request from a client computer to access a document stored at a particular URL. Browser software such as Netscape.TM. Navigator, of Netscape Communications Corporation of Mountain View, Calif., is well known and widely available. When the server receives the user's request, it sends the requested HTML Web document to the client where the document can be displayed. The communications protocol used in making such a request and in transferring Web documents is \"Hypertext Transfer Protocol\" (HTTP). For more information about the Web, see for example T. Berners-Lee, R. Cailliau, A. Loutonen, H. F. Nielsen, and A. Secret, \"The World Wide Web,\" Communications of the ACM, vol. 37(8), August 1994.\nAs a theoretical example of a Web site, consider a particular business, Business A, having a Web site dedicated expressly to providing information about Business A. Business A's Web site may contain an introductory Web page describing the general technologies and product areas in which Business A is involved. The introductory Web page may also contain links to other Web pages related to Business A. For example, the introductory Web page may have a link to a separate Web page disclosing financial aspects, such as the annual report, stock price history, and the like, of Business A. By selecting the link, the linked Web page is sent to the client computer and displayed to the end user. Thus, a user initially contacting Business A's introductory Web page may discover various other information about Business A by virtue of the linked nature of the Web.\nAlthough a linked Web document is probably related to the original Web document, the linked Web document may reside separately from the original Web document. That is, the linked Web document may reside in a directory which is physically separate from the directory in which the original Web document is stored. In fact, the linked Web document may even reside in a computer or in hardware separate from the computer or hardware in which the original Web document is stored. Therefore, multiple Web documents may be linked even when the Web documents are not stored contiguously or even proximately. Separate storage of linked Web documents may result from the fact that a particular storage device does not have the capacity to contain all of the linked documents, or from the fact that a collection of linked web documents are maintained and/or owned by another individual and/or company.\nAlthough on-line Web sites have substantial utility as information providers in both personal and professional environments, on-line access is not available to many potential end-users. For example, in government or high security environments, computing terminals are \"quarantined\" or isolated from the outside world. In such instances, on-line access is not available or is prohibited. Additionally, on-line access is often simply not feasible for any one of numerous reasons. For example, an end-user's hardware/computer system may not be adapted for on-line use, or may operate too slowly to effectively convey information to the end-user over an on-line connection. Likewise, if an end-user's client computer is located very distantly from the server system, on-line access may be prohibitively slow. As yet another example, during peak Internet or on-line use periods, response from the server system may be so greatly delayed as to render on-line retrieval of information impractical.\nThus, a need exists for a system which allows end-users to access Web site information, or other types of information commonly presented in an on-line manner, when the end-user is off-line. A further need exists for a system in which an off-line end-user can access linked information even when the linked information is so voluminous as to necessitate storage on multiple storage devices. Still another need exists for a system in which an off-line end-user accesses information in a manner which emulates on-line retrieval of the information."}
{"text": "P1 phosphor is the dielectric material commonly used as a viewable bistable storage target for storage cathode ray tubes such as taught in U.S. Pat. Nos. 3,293,433; 3,293;474; 3,531,675 and U.S. patent application Ser. No. 356,029, filed Apr. 30, 1973. This phosphor dielectric material has been found to have desirable dielectric characteristics to enable information written thereon via a deflected cathode ray beam of electrons to be stored for any length of time until it is desired to erase such stored information. This P1 phosphor material can also be used to display information via the deflected beam of electrons in a non-store mode of operation. Hence, the phosphor material provides bistable operation in store or non-store modes of operation depending upon which mode of operation is desired, and brightness of the visible image formed on the phosphor layer as well as image contrast, writing speed and erase speed are important features attributed to P1 phosphor.\nWhile P1 phosphor has provided excellent operating characteristics as a bistable storage target for cathode ray tubes, the operating life of such target is not as great as desired. Operating life of CRT targets are dependent upon operator use and uniformity of screen usage.\nU.S. Pat. No. 2,241,939 discloses manganese-activated zinc orthosilicate (P1) phosphor that is mixed with metal oxides, some of which can be rare earth oxides, which mixture is then subjected to heat thereby forming a specialized sensitized phosphor material for use in a mercury lamp to provide higher luminance efficiency.\nU.S. Pat. No. 3,862,450 teaches the concept of bonding or adhering secondary electron emissive material on the surface of phosphor storage dielectric particles that are used to form the storage dielectric phosphor layer of a storage target of a direct-viewing bistable storage tube. Such storage dielectric layer increases the writing speed of the storage target and also the life thereof, but the increased target life is not acceptable.\nU.S. patent application Ser. No. 658,977, filed Feb. 18, 1976 discloses the use of an admixture of P1 phosphor particles and yttrium oxide, or yttrium oxysulfide, or rare earth activated yttrium oxide or yttrium oxysulfide to provide a bistable storage target of greater life than that of the storage target of pure P1 phosphor; however the target life is not as long as directed."}
{"text": "1. Field of the Invention\nThe invention relates generally to systems and methods to deliver crew, cargo and other logistics services to space platforms. More particularly, the invention relates to systems and methods for automated rendezvous and docking of a spacecraft to a space platform.\n2. Related Art\nThe Space Shuttle has provided a majority of the assembly and logistics service support to the International Space Station (ISS). The shuttles are expensive to operate, and they place personnel at risk even when a primary payload does not otherwise require a crew. The Russian Soyuz and Progress spacecraft are also used for crew and cargo transport, respectively, to the ISS, including the transfer of consumables such as propellant. Both of these systems have a lengthy heritage and have driven some of the design of the ISS they now support. The European and Japanese space agency transfer vehicles currently in development will be capable of transporting cargo to and from the ISS as well. The development of these systems has helped to identify opportunities to improve on-orbit operations.\nA spacecraft is designed to fulfill a set of requirements which reflect mission objectives as well as many constraints. The spacecraft systems are designed to carry out the mission and support the needed cargo. While there are general requirements, detailed requirements are honed to meet specific mission requirements and other considerations are made with regard to performance capabilities, flight environments, safety, mechanical and electrical interfaces, guidelines, standards, regulations, flight rules, schedule, and budget, for example. A number of alternatives may meet the requirements and the systems and subsystems are designed and sized to best satisfy the requirements. The resulting structures and components are modeled, simulated and constructed. Each of these individually, and integrated as a system, must pass through various verification, validation, qualification, certification and acceptance tests. The process is typically long and expensive. The complexity is compounded by the number of systems, their interactions with each other and with other systems and spacecraft. Even the smallest of changes at an inopportune point in the process is potentially very expensive. That being said, large sums are spent to save even larger sums or to mitigate risks.\nThe cost of spacecraft and their missions is unavoidably tied to the delivering mass to orbit. So, the mass of the spacecraft and its various components is of great concern. The spacecraft and their cargo become high-value assets, and the risks to these mitigated to the greatest degree possible by various logistical measures which, in turn, make the systems and procedures even more expensive to undertake. The result is that simplifications to spacecraft and their missions can greatly affect (reduce) their cost.\nOne such simplification is to relieve a spacecraft of some part of its functionality or mission, for example, by placing that responsibility with another system or spacecraft. If the other system or spacecraft is reusable, the expense associated with that part of functionality or mission can be amortized over a number of spacecraft or missions; on-orbit operations is one area in which this can be done. The orbital maneuvering, and automated rendezvous and docking (AR&D) functions are candidates for incorporation into a reusable spacecraft.\nTransfer vehicles such as the Orbital Maneuvering Vehicle (OMV) have been proposed. The OMV is used to change the orbital plane of a second spacecraft, or to boost the second spacecraft into a higher orbit. This allows the second spacecraft to be built with much smaller propulsive systems designed primarily for attitude control (reaction control systems). The OMV stays in orbit and is able to provide its services to several spacecraft. The result is that the engines used for the orbital changes need only be launched into orbit once, not with each spacecraft. The OMV must carry enough propellant for all of its missions, but the systems associated with carrying the propellant (the “overhead”) are amortized over all of the missions. Also, additional reserve propellant is usually carried by a spacecraft to enable the spacecraft to handle a worst-case scenario. That reserve is likely never used so the cost of putting it into orbit is “wasted”. The OMV is able to use that reserve propellant for another mission.\nAttempts are being made to simplify the process of mating visiting vehicles such as logistics spacecraft to the ISS. One such attempt involves the free-flying capture of a visiting vehicle with a subsequent berthing rather than relying on the visiting vehicle to dock directly with the ISS. When docking, a visiting vehicle must safely align itself with a docking port and drive into it with enough force to properly mate with the docking mechanism but not so much as to damage the ISS or force it “out of control”. A berthing operation requires the visiting vehicle to rendezvous with the ISS and to move to within reach of the ISS's robotic arm subject to a number of constraints. The robotic arm then “grabs” and pulls the visiting vehicle to a docking port. The free-flying capture and berthing procedure, however, has proven to be more problematic than hoped and better means for docking visiting vehicles are being explored.\nIn the early 1990's, a tug was proposed for transferring a supply spacecraft from its orbit to the proximity of the International Space Station (ISS). The features of the proposed system did not justify its manufacture and it was never built. A tug travels to the supply spacecraft and grabs or mates with it. The tug then “pushes” the supply spacecraft to the ISS. The supply spacecraft is then docked. The supply spacecraft is driven into a compatible and available docking port where the supply spacecraft's docking mechanism mates directly with the ISS. The tug is then free to go or may remain to remove the supply spacecraft.\nThe visiting vehicle must have an ISS compatible docking mechanism and also have appropriate fixtures for the tug to attach to the supply spacecraft. In the case of a supply vehicle needing to perform a fuel transfer, the appropriate ISS port would be a port with a Common Berthing Mechanism (CBM) or the Russian Probe and Cone Mechanism (RPCM).\nThe ISS is also equipped with an Androgynous Peripheral Attachment System (APAS) which is used by the Shuttles. The ISS is equipped with 2 such mechanisms. It is notable that the both the CBM and APAS are no longer manufactured and that a limited number are available to be used on supply spacecraft. The preservation of these mechanisms is highly desirable. After resupply, the current roster of supply spacecraft deorbit and burn up in the atmosphere (e.g., HTV, ATV, Progress), hence the mechanism for docking with the ISS is lost. A Crew Exploration Vehicle would deorbit and be recovered as are other manned spacecraft such as the Shuttle or Soyuz. However, the deorbit operations may require the docking mechanism be jettisoned."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an information processing apparatus, an information processing system, and a data conversion method. The present invention also relates to a non-transitory computer-readable storage medium having stored therein a program for causing a computer to execute an information processing or a data conversion process.\n2. Description of the Related Art\nConventionally, electronic data in a certain data format may be converted into electronic data in another data format, such as when a printer driver converts an office document into print data, for example. A data conversion process to convert the electronic data into a format that is suited for use at the time of outputting the electronic data is known from a Japanese Laid-Open Patent Publication No. 2006-81081, for example.\nIn addition, the data conversion process is not limited to that performed in a single computer, and the data conversion process may be performed in a computer system that uses a plurality of computers, as proposed in a Japanese Laid-Open Patent Publication No. 2009-75689, for example.\nWhen making the conversion from one data format into another data format, a data conversion engine (or a data conversion program or a data conversion library) may be used in order to improve the independence or maintenance ease of the program.\nHowever, the number of types of data conversion engines that are required to perform this conversion corresponds to the number of types of pairs formed by the number of types of input data formats and the number of types of output data formats. For example, when there are three (3) types of data format candidates for the data format before the conversion and there are three (3) types of data format candidates for the data format after the conversion, nine (9) types of data conversion engines are required in accordance with the number of types of pairs formed by the number of types of input data formats and the number of types of output data formats. In other words, the larger the number of options for the data formats before or after the conversion, the larger the number of types of data conversion engines that need to be prepared.\nHence, as the number of options for the data formats before or after the conversion increases, the number of types of data conversion engines that need to be prepared increases. As a result, when the number of options for the data formats before or after the conversion is relatively large, it may be difficult to flexibly cope with the increase in the number of types of data conversion engines that need to be prepared."}
{"text": "Human enterovirus D68 (EV-D68) was first isolated from samples obtained in California in 1962 from four children with pneumonia and bronchiolitis. The type strain isolated from one of these children has been designated the Fermon strain. Subsequently, only small numbers of EV-D68 cases were reported until the early 2000s. However, from 2008-12 outbreaks in Japan, the Philippines, the Netherlands, and the USA (Georgia, Pennsylvania, and Arizona) have revealed EV-D68 as an emerging pathogen capable of causing severe respiratory illness. During the 2014 enterovirus/rhinovirus season in the United States, EV-D68 circulated at an unprecedented level. From August 2014 to January 2015, CDC and state public health laboratories confirmed a total of 1,153 cases of respiratory illness caused by EV-D68, with at least 14 deaths. Infected individuals were primarily children, and resided in 49 states and the District of Columbia. The CDC has also reported there were likely millions of EV-D68 infections in which the etiology was not determined.\nIn mid-August of 2014, hospitals in Missouri and Illinois noticed an increased number of patients with severe respiratory illness and reported the presence of EV-D68. Because efforts to define the outbreak were hampered by the lack of a test for EV-D68 that did not require nucleotide sequencing, there is a need in the art for a rapid, specific RT-PCR assay."}
{"text": "Field\nThe following description relates to an apparatus and transmission method for transmitting communication signals at the maximum data transmission rate within a range to satisfy the specific absorption rate (SAR) regulations associated with electromagnetic wave exposure from a portable device.\nDiscussions of Background\nWith the development of data communication and display technologies, portable devices (for example, smartphones and tablet PCs) which enable data communication have become a necessity for modern people. Such portable devices operate based on electronic components and thus expose users to electromagnetic waves. Recently, as the harmfulness of electromagnetic waves to the human body has become an issue, specific absorption rate (SAR) regulations for portable devices have been established. SAR regulations use energy absorbed per unit time and mass as an index, and regulate the rate of absorption of electromagnetic wave energy by a biological tissue, and is measured in W/kg. In the Republic of Korea, the limit SAR value for portable devices is set at 1.6 W/kg, and since April 2002, a SAR limit value for the human body has been introduced to protect the human body against electromagnetic waves absorption. Therefore, SAR values of portable devices have been required to be tested and the sale of portable devices with a SAR value higher than the limit value has been prohibited. In the Republic of Korea, the government has further proposed amendments to heighten the SAR regulations. In the amended SAR regulations, Body SAR regulations which require measuring SARs of the six sides of a portable device have been added, and a grading system has also been implemented.\nAs illustrated in FIG. 1, a circuit switched fall-back (CSFB) type portable device is designed to perform communication only in a single frequency band through a main antenna 30 even when two different frequency bands are supported. For example, the portable device is designed such that in a case where transmission of a voice signal 11 is requested while transmission of a data signal 12 is performed, a channel regarding the transmission of the data signal 12 is disconnected by a switch 20, a channel regarding the transmission of the voice signal 11 is connected by the switch 20, and the main antenna 30 performs communication only in a single frequency band. Since the existing SAR regulations only required SAR measurement at a human's head, for a CSFB type portable device, the SAR regulations could be easily satisfied by positioning the main antenna 30 at a position farthest from the head.\nHowever, the CSFB type portable device has a disadvantage that communication in a certain frequency band may be interrupted. In the above example, transmission of voice signal 11 is interrupted by communication of the data signal 12. Therefore, the simultaneous voice and LTE (SVLTE) technology which simultaneously enables transmission and reception in two or more frequency bands has been introduced. However, since the SVLTE technology uses a CDMA network and a LTE network at the same time, there is a limit in improving data transmission rates. Accordingly, the SVLTE technology has been further developed into LTE carrier aggregation (CA) technology which performs transmission and reception simultaneously using two LTE networks.\nReferring to FIG. 2, a portable device which uses the LTE CA technology has a diplexer 20 inserted instead of a switch as shown in FIG. 1. The diplexer 20 filters each signal, and may transmit two LTE signals 11, 12 in different frequency bands by simultaneously filtering the two LTE signals. Accordingly, voice communication and data communication can be simultaneously performed over the LTE network, and thus data signals can be continuously received without interruption even during voice communication.\nSince strength of electromagnetic waves increase as the power of transmission signals increase, in order for the LTE CA type portable device to comply with the SAR regulations, the power of a signal associated with a certain frequency band may need to be limited. For example, if the transmission power of the voice signal 11 is equal to or higher than a predetermined magnitude (for example, 20 dBm or higher), even if a base station commands the portable device to transmit the data signal 12 at its maximum magnitude (for example, 23 dBm), the portable device may need to control the magnitude of the transmission power of the data signal 12 to be set at a predetermined magnitude (for example, 19 dBm) or less.\nThe LTE CA type portable device has problems that in a poor electromagnetic wave environment, the data transmission rate may decrease, the signal quality of data communication may be degraded, and/or the communication may be interrupted. To resolve these problems, considering the fact that the electromagnetic wave absorption rate can be decreased by increasing the distance between the transmission antennas, a transmission method using a plurality of transmission antennas may be considered. However, since multipurpose antennas such as a GPS, Bluetooth, and Wi-Fi all need to be integrated in the portable device, there is a limit in satisfying the SAR regulations by increasing the distance between the multiple antennas."}
{"text": "1. Field of the Invention\nThis invention relates to an image processing system which outputs output image data having an output gradient number of n lower than the input gradient number m by carrying out multi-valued processing using a dither matrix on input image data having an input gradient number of m.\n2. Description of the Related Art\nWhen image data representing information of an image created by a personal computer or image data obtained by reading an original with a scanner or the like is to be output to a printer to be printed thereby or to a display to be displayed thereby, processing to conform the image data to the lower number of gradation of the output side and to have the same gradation as the information of the image or the original in appearance is generally carried out on the image data since the number of gradation of the output side device, the printer or the display, is lower than that of the information of the image created by the personal computer or the image data read-out by the scanner.\nFor example, when the multi-valued image data created by the personal computer is to be printed by a two-valued printer, an organizational dither method is employed. The organizational dither method is a method in which n×n pixels are considered to be the smallest unit of expressing the gradation, and a n×n matrix of threshold values (will be referred to as “a dither matrix”, hereinbelow) corresponding thereto is prepared, and the dither matrix is superposed on the image data as a kind of mask, whereby the density of each pixel is compared with the corresponding threshold value and “1” is output when the former is larger than the latter while “0” is output when the former is smaller than the latter. Then after the same processing is finished on a group of n×n pixels, the dither matrix is moved to another group of n×n pixels and the same processing is repeated, whereby each components of the whole image data is turned to one of the two values, “1” or “0”.\nFurther, when each pixel is turned to one of the multi-values (multi-valued) in place of the two values, there is proposed, in U.S. Pat. No. 5,640,249 (will be referred to as “the first paper”, hereinbelow.), a method where the density of each pixel is compared with a plurality of threshold values. For example, a plurality of dither matrices each having a threshold value necessary to the multi-valued processing (turning the density of each pixel to the multi-values) are prepared from one basic dither matrix, and the density of each pixel is multi-valued by comparing the density of each pixel with a plurality of threshold values on the same coordinates in a plurality of dither matrices. That is, in accordance with the first paper, (x−1) comparisons are necessary when the number of gradations into which each pixel is to be multi-valued is x.\nWhereas, in “Digital Color Halftoning” by Henry R. Kang., IEEE PRESS, SPIE ISBNO-8194-3318-7, pp. 445-449, 1999/11/1 (will be referred to as “the second paper”, hereinbelow.), there is disposed a method where multi-valued processing can be realized by one comparison. More particularly, the method disclosed in the second paper is a method of obtaining output image data O(x,y) by the use of the following formula (3). Since the multi-valued processing can be executed according to one formula without using comparison operation, the method of the second paper can be executed at a higher speed as compared with the method of the first paper. Further, the method of the second paper requires less memory capacity for storing the dither matrix as compared with the method of the first paper where a plurality of dither matrices are used.O(x,y)=int[α+d(xd,yd)/(gmax+1)]  (3)wherein α=(n−1)/(m−1)*I(x,y), n represents output gradient number, m represents input gradient number, O(x,y) represents output image data (0≦O(x,y)≦(n−1)), I(x,y) represents input image data (0≦I(x,y)≦(m−1)), gmax represents the maximum value of the dither matrix, d(xd,yd) represents the value of the dither matrix (0≦d(xd,yd)≦gmax), xd=x mod dsize_x,    yd=y mod dsize_y, dsize_x represents the size of the dither matrix in x direction, dsize_y represents the size of the dither matrix in y direction, and int[ ] is an integer-valued function.\nHowever, if the method disclosed in the first paper is employed, there arises a problem that, for instance, when a gradation image is to be output, the image data components at a particular density can be processed to one gradation and output in the gradation though the image data components at almost all the densities can be processed to two gradations and output in the gradations. The result obtained by carrying out the multi-valued processing on a gradation image by the use of, for instance, the dot dispersion type dither matrix and the above formula is shown in FIG. 5. As shown in FIG. 5, though, in a lateral view, almost all the regions are expressed in two different gradations, a region which is expressed in a single gradation exists near the center. When such an image is output, e.g., printed, though almost all the regions have a feeling of granulated material, the region expressed in a single gradation has no feeling of granulated material, whereby people visually feel as if there is different material in the region, and will feel unnatural."}
{"text": "1. Field of the Invention\nThe present invention relates to method for driving a flat panel displays, and particularly to a method for driving a flat panel display for reducing power consumption of a source driver by modified scan signals, and therefore for reducing the system power consumption.\n2. Description of Related Art\nIn order to catch up with the modern lifestyle, video and image devices are becoming slimmer and lighter. Despite of advantages it may have, a conventional cathode-ray tube (CRT) display not only is large in bulk that occupies too much room because of its intrinsic structure of the electronic cavity, but also radiates rays which may hurt human eyes. Therefore, accompanying the development of optoelectronic technology and semi-conductor processing technology, flat panel displays including liquid crystal display (LCD), organic light-emitting diode (OLED) display and plasma display panel (PDP) are gradually becoming a mainstream in the display market.\nResolutions and refreshing frequencies of flat panel displays are continuously improving. Consequently, refreshing frequencies of scan lines are demanded to be more and more rapid, which contradicts the designs for power saving by system engineers. As a result, a technology for eliminating the contradiction therebetween called “smart charge sharing” technology is developed thereby.\nFIGS. 1A and 1B are schematic diagrams of a conventional charge sharing technology provided by Dr. Craig Zajac of the US National Semiconductor Incorporation. Referring to FIG. 1A, first, a flat panel display 100 includes a source driver 110 and a pixel array 130. The pixel array 130 includes a plurality of data lines DL0 through DLm electrically connected to the source driver 110.\nThe source driver 110 includes a plurality of output amplifiers A1 through Am. Each of the output terminals of the amplifiers is respectively connected to a corresponding data line. The source driver 110 further includes a plurality of switches SW1 through SWm−1 for connecting adjacent two data lines. For example, a switch SW1 is adapted for connecting the adjacent data lines DL0 and DL1. As shown in FIG. 1A, each data line is taken as a sum of loads of resistance and capacitance of a corresponding output amplifier.\nBefore the source driver 110 driving the pixel array 130, voltages of each pair of the adjacent data lines are respectively higher or lower than a common voltage. Meanwhile, all of the switches SW1 through SWm−1 are at turn-off status. At the instance that the source driver 110 starts to drive the pixel array 130, all of the switches SW1 through SWm−1 will be switched to turn-on status as shown in FIG. 1B. At this certain instance, all amplifiers A1 through Am are at disable status, which is called Hi-Z mode during which the output amplifiers A1˜Am without current consumption. However, as the switches SW1 through SWm−1 are at turn-on status, there will be a current flowing from a data line having a voltage higher than the common voltage to a data line having a voltage lower than the common voltage, the path of which is as illustrated of the arrowheads in FIG. 1B. Thus, charges can be neutralized therein. After experiencing a Hi-Z mode, the switches SW1 through SWm−1 will go back to the turn-off status, while the source driver 110 can the drive pixel array 130 as usual.\nIn summary, the principle of the charge sharing technology is to reallocate energy (charges) stored in the data lines and whereby to drive the scan lines to a half of the final value without power consumption. Such technology is enough for those small flat panel displays taking little loads and providing a lower resolution. Unfortunately, whenever an overall slew rate of the source driver 110 slows down, an efficiency of the source driver 110 will drop down accordingly. Therefore, such a source driver can not satisfy the requirements of the new generation flat panel displays which demand higher resolution, higher load and larger size."}
{"text": "In recent years, a higher integration of a semiconductor device requires the narrower wiring and the multilayer wiring, and hence it is necessary to make a surface of a semiconductor substrate highly planarized. Specifically, finer interconnections in highly integrated semiconductor devices have led to the use of light with shorter wavelengths in photolithography, so that a tolerable difference of elevation at the focal point on the substrate becomes smaller in the light with shorter wavelengths. Therefore, a difference of elevation at the focal point should be as small as possible, i.e., the surface of the semiconductor substrate is required to be highly planarized. One customary way of planarizing the surface of the semiconductor substrate is to remove irregularities (concaves and convexes) on the surface of the semiconductor substrate by a chemical mechanical polishing (CMP) process.\nIn the chemical mechanical polishing process, after a surface of a semiconductor substrate has been polished for a certain period of time, the polishing process should be finished at a desired position or timing. For example, some integrated circuit designs require an insulating film (layer) of SiO2 or the like to be left on a metallic interconnection of copper, aluminum, or the like. Since a metallic layer or other layers are further deposited on the insulating layer in the subsequent process, such an insulating layer referred to as an interlayer. In this case, if the semiconductor substrate is excessively polished, the lower metallic layer is exposed on the polished surface. Therefore, the polishing process needs to be finished in such a state that a predetermined thickness of the interlayer remains unpolished.\nAccording to another polishing process, interconnection grooves having a certain pattern are formed in a surface of a semiconductor substrate, and after a Cu layer is deposited on the semiconductor substrate, the interconnection grooves are filled with copper or copper alloy, and then unnecessary portions of the Cu layer are removed by a chemical mechanical polishing (CMP) process. Specifically, the Cu layer on the semiconductor substrate is selectively removed by the chemical mechanical polishing process, leaving only the Cu layer in the interconnection grooves. More Specifically, the Cu layer is required to be removed until an insulating layer of SiO2 or the like is exposed in surfaces other than the interconnection grooves.\nFurther, in some cases, interconnection grooves for a predetermined wiring pattern are formed in a semiconductor substrate, conductive materials such as copper (Cu) or copper alloy are filled in such grooves of the semiconductor substrate, and then unnecessary portions of the conductive materials on the surface of the semiconductor substrate are removed by a chemical mechanical polishing (CMP) process. When the copper layer is polished by the CMP process, it is necessary that the copper layer on the semiconductor substrate be selectively removed therefrom in such a state that the copper layer is left in the grooves for a wiring circuit, i.e. the interconnection grooves. More specifically, the copper layer on those surfaces of the semiconductor substrate other than the interconnection grooves needs to be removed until an insulating layer of SiO2 or the like is exposed on the polished surface.\nIn such cases, if the semiconductor substrate is excessively polished until the Cu layer in the interconnection grooves is removed together with the insulating layer, then the resistance of the circuits on the semiconductor substrate would be so increased that the semiconductor substrate might possibly need to be discarded, resulting in a large loss of resources. Conversely, if the semiconductor substrate is insufficiently polished to leave the copper layer on the insulating layer, then interconnections on the semiconductor substrate would not be separated from each other as desired, but a short circuit would be caused between those interconnections. As a result, the semiconductor substrate would be required to be polished again, and hence its manufacturing cost would be increased. The above problems also occur when another metallic film of aluminum or the like is formed on a semiconductor substrate and polished by the CMP process.\nTherefore, it has been proposed to detect an end point of the CMP process with use of an optical sensor. Specifically, an optical sensor comprising a light-emitting element and a light-detecting element is provided in a polishing apparatus. The light-emitting element of the optical sensor applies light to the polished surface of a semiconductor substrate, and the light-detecting element detects a change in reflectance of the light reflected from the polished surface for thereby measuring the thickness of an insulating layer or a metallic layer on the polished surface. Thus, the end point of the CMP process is detected from the measured film thickness.\nIn a polishing apparatus for performing the CMP process, a polishing pad mounted on the upper surface of a polishing table usually has a low light transparent. Therefore, in the case where light from the optical sensor is to be applied from below the polishing pad to the polished surface of a semiconductor substrate placed on the polishing table, a light-transmittable window having a high light transmittance for allowing light to pass therethrough is provided in the polishing pad, and light from the optical sensor is applied through the light-transmittable window to the polished surface of the semiconductor substrate.\nFIG. 1 is an enlarged fragmentary cross-sectional view showing a conventional double-layer polishing pad including a light-transmittable window. As shown in FIG. 1, a polishing pad 310 comprises an upper-layer pad 311 and a lower-layer pad 312. The upper-layer pad 311 has a hole 311a defined therein, and a light-transmittable window 341 is disposed in the hole 311a of the upper-layer pad 311. The lower-layer pad 312 has a light passage hole 312a defined therein. The light passage hole 312a has a diameter smaller than the hole 311a. The holes 311a, 312a having different diameters provide a step 313 therebetween.\nThe conventional polishing pad 310 is manufactured as follows: An adhesive agent is applied to the lower surface of the upper-layer pad 311 and the upper surface of the lower-layer pad 312, and the upper-layer pad 311 and the lower-layer pad 312 are vertically pressed against each other. Thus, the upper-layer pad 311 and the lower-layer pad 312 are bonded to each other. Then, the light-transmittable window 341 is fitted into the hole 311a. The light-transmittable window 341 is bonded to the lower-layer pad 312 by the adhesive agent applied to the upper surface of the step 313.\nHowever, since the light-transmittable window 341 is bonded to the lower-layer pad 312 only at the upper surface of the step 313, its bonded area and hence bonding strength is small. Therefore, the light-transmittable window 341 may possibly be peeled off the polishing pad 310. Depending on the forces applied to the polishing pad 310, the light-transmittable window 341 may not be peeled off in its entirety, but may partly be peeled off and a gap is produced between the light-transmittable window 341 and the lower-layer pad 312 or the step 313. The gap thus formed permits a polishing liquid used on the upper surface of the polishing pad 310 to leak onto the lower surface of the light-transmittable window 341. When the polishing liquid is attached to the lower surface of the light-transmittable window 341, the reflectance of the light-transmittable window 341 is greatly reduced, thereby making it difficult to detect a change in reflectance of the polished surface of the semiconductor substrate with the optical sensor. In this case, the film thickness of the semiconductor substrate cannot be measured with high accuracy.\nIf the polishing liquid is introduced into the gap between the light-transmittable window 341 and the step 313, then a region where the elasticity is not uniform is produced in the polishing pad 310 due to the swelling thereof. This region may adversely affect the polishing process of the semiconductor substrate. Furthermore, if the polishing liquid having a low transparent is irregularly introduced into between the light-transmittable window 341 and the optical sensor, then signals detected by the optical sensor become unstable, making the detected result less reliable.\nAs described above, the light-transmittable window 341 is fitted in the hole 311a defined in the upper-layer pad 311. The hole 311a has dimensions slightly greater than the light-transmittable window 341 to allow the light-transmittable window 341 to easily be fitted in the hole 311a. Therefore, after the light-transmittable window 341 has been placed in the hole 311a, there is a small gap 314 between the light-transmittable window 341 and the hole 311a. Accordingly, the polishing liquid tends to be introduced into the small gap 314 and become solidified in the small gap 314. The solidified polishing liquid may cause scratches on the semiconductor substrate which is polished on the upper surface of the polishing pad 310."}
{"text": "1. Field of the Invention\nThe present invention relates to a negative resist composition and a patterning process using the same.\n2. Description of the Related Art\nIn recent years, as an integrated circuit progresses toward a high integration, a further finer patterning is required. Especially, when a resist pattern of 0.2 μm or less is formed, as a usual manner, a chemically amplified resist composition wherein an acid generated by an exposure light acts as a catalyst, from which a high sensitivity and a high resolution can be obtained, has been used. As the exposure light source, a high energy beam such as an ultraviolet beam, a far-ultraviolet beam and EB has been used, and especially, an exposure method by using EB or EUV has been drawing attention as the method that the finest pattern is expected among the methods currently under consideration.\nIn a resist composition, there are a positive type in which an exposed part is dissolved and a negative type in which an exposed part remains as a pattern. They are selected based on the ease of its use and depending on the required resist pattern. A chemically amplified negative resist composition usually contains a polymer which is alkaline-soluble and is made alkaline-insoluble by action of an acid, and an acid generator generating an acid by decomposition with an exposure light, and, in addition, usually contains a basic compound to control diffusion of the acid generated by the light exposure.\nMany negative resist compositions using a phenolic unit as an alkaline-soluble unit constituting the polymer that is soluble in an aqueous alkaline developer have been developed, especially for exposure to a KrF excimer laser beam. These have not been used for an ArF excimer laser beam because a phenolic unit does not transmit the light when the exposure light having a wavelength of 150 to 220 nm is used. However, in recent years, these have been drawing attention again as a negative resist composition for exposure to EB and EUV, which are exposure methods to obtain a further finer pattern, as reported in Japanese Patent Laid-Open (kokai) No. 2006-201532, Japanese Patent Laid-Open (kokai) No. 2006-215180, and Japanese Patent Laid-Open (kokai) No. 2008-249762.\nAs a further finer pattern is required, negative resist compositions using a hydroxystyrene unit which is the typical structure of the phenolic unit have been improved many times. However, as a pattern size is miniaturized considerably like 0.1 μm or less, conventional materials cannot have solved the problems that a so-called bridge (a finely-filamented resist layer remained between pattern lines of a fine pattern) is easily caused and an undercut occurs at the substrate's interface and thus patterns easily collapse.\nIn addition, as for so-called substrate dependence of a pattern (a pattern profile change near the substrate due to a material of a substrate to be processed), with a miniaturization of a pattern to be obtained, even a small profile change has become a problem. Especially, when pattern forming is done by using a chemically amplified negative resist on a chromium oxynitride, the outermost surface material of a photomask blank, in processing of the photomask blank, a so-called undercut (a cleavage formed at the contact part of the resist pattern with the substrate) is caused, but conventional materials could not solve it sufficiently.\nFurther, in the case of the resist elaboration as mentioned above, examples of characteristics to be demanded for resist compositions are not only a higher resolution which is a fundamental performance of a resist, but also a higher etching resistance. This is because a progressive miniaturization of patterns necessitates to progressively decrease thicknesses of resist films. Known as one technique for obtaining such a higher etching resistance, is a method, which also has been disclosed in Japanese Patent Laid-open (Kokai) No. 2008-249762, to introduce a compound: a polycyclic compound such as indene and acenaphthylene, which includes an aromatic ring(s) and a non-aromatic ring(s) having a carbon-carbon double bond conjugated with the aromatic ring, into a polymer having a hydroxystyrene unit as a subsidiary component.\nIn turn, as polymers for positive resists, it has been proposed to use a polymer having only an indene skeleton such as described in Japanese Patent Application Laid-open (kokai) No. 2004-149756, while another method has been proposed in Japanese Patent Application Laid-open (kokai) No. 2006-169302 to use a unit having an acenaphthylene skeleton by combining it with a hydroxystyrene derivative.\nIn Japanese Patent Laid-Open (kokai) No. 2008-249762, a polymer comprising a hydroxystyrene unit and a chlorostyrene unit was used as a styrene derivative unit, while a miscible system of triphenyl sulfonium-2,5-dimethylbenzene sulfonic acid and triphenyl sulfonium-2,4,6-triisopropylbenzene sulfonic acid was used as an acid generator. By using this material containing a crosslinker which forms cross-linkages among polymers by the aid of the acid as a catalyst to make the polymer insoluble in a developer, a line-and-space pattern of 60 nm with few undercuts and without a bridge has been obtained. However, pattern formation of 50 nm or less was difficult because of occurrence of a bridge."}
{"text": "Unless otherwise indicated herein, the approaches described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nFIG. 1 depicts a reference clock doubler circuit 100 and FIG. 2 shows waveforms of the reference clock doubler circuit 100. A buffer 102 outputs a full-swing clock, C1. A delay element 104 delays the full-swing clock C1 and outputs a delayed full-swing clock C2b. A logic element 106 (e.g., XOR gate) outputs a doubled reference clock C2x. The rising edges of the doubled reference clock C2x comes from the rising and falling edges of the C1 clock. The falling edges of the doubled reference clock C2x come from the rising and falling edges of the delayed full-swing clock C2b. The non-50% duty cycle in clock signal C1 causes a mismatch between the even clock period, Tref, even, and odd clock period, Tref, odd, in the doubled reference clock C2x. For example, the mismatch causes high-reference spur. Additionally, a higher average “on time” is forced for charge pumps in a phase lock loop (PLL) using the doubled reference clock signal and also for time-to-digital converters (TDCs) for digital PLLs, which increases the phase noise due to the charge pump or TDC thermal and flicker noise.\nFIG. 3 depicts a PLL 109 illustrating the reference doubler problem. A phase frequency detector (PFD) 110 receives a reference clock clkref (e.g., doubled reference clock C2x) and a feedback clock clkfb. Phase frequency detector 110 outputs an up and down signal that represents the difference in phase between the two input signals. An up and down signal is output to a charge pump 112. Charge pump 112 charges capacitors C1 and C2 of a loop filter when an up switch 506-1 is closed and dissipates charge when a down switch 506-2 is closed. VCO 114 receives a tuning voltage from the loop filter and generates an oscillating signal with a frequency controlled by tuning voltage. A divider 116 divides the output of VCO 114 to generate the feedback signal clkfb.\nAt 117, a summary of the waveforms for signals received and output at PFD 110 is shown. At 118, reference clock clkref is shown; at 120, feedback clock clkfb is shown; at 122, an up signal is shown; and at 124, a down signal is shown. The even and odd cycles of reference clock clkref are not the same length. The up signal is generated by PFD 110 when a rising edge of reference clock clkref is before a rising edge of the feedback clock clkfb. The down signal is generated when the rising edge of the reference clock clkref is after a rising edge of the feedback clock clkfb. The up and down signals alternate every even and odd clock period.\nThe up and down signals deliver charge to capacitors C1 and C2 that cancel each other. However, uncorrelated noise may be received and increases in-band noise. For example, at 128, VCO 114 up-converts the noise received into reference spurs. One solution to correct the mismatch is to add digital and analog tuning to adjust the duty cycle of reference clock clkref. However, the analog tuning increases the noise. Also, digital tuning cannot track temperature variation without introducing large glitches to the output clock."}
{"text": "This invention relates to a machine for automatically retrieving and dispensing a selected container from a store of containers and for receiving and storing a container which is returned to the machine. The invention is particularly useful as a vending machine for video cassettes and the like which are dispensed for the temporary usage of a person and which must be returned and stored in the machine for subsequent dispensing.\nIn the recent years the increased usage of video cassette recorders (VCR's) to play prerecorded programming, primarily in the form of popular movies, has resulted in a large rental market for such video cassettes. The distribution chain for video cassette rentals has primarily been through a plethora of video cassette rental stores which have opened all over the country to meet this demand. A typical transaction involves going into one of such video cassette rental stores, engaging a sales person to orally request a selected movie, often after waiting in a long line, and if the selected video cassette is in stock completing a rental transaction. Upon returning the rented video cassette to the store often times it is again necessary to engage a sales person, possibly after waiting in another line, to settle the rental account. In short, the principal mechanism which has developed in the market place for the rental of video cassettes is often time consuming and inconvenient.\nBecause of its relatively small size, the standard video cassette is capable of being disposed by a vending machine. In fact, applicant is aware of one such vending machine which has been developed for this purpose. This known video cassette vending machine is basically a cigarette-type vending machine which has been modified to accommodate video cassettes. This machine contains a limited store of cassette movie titles which are selected and dispensed in a manner similar to the selection and dispensing of a package of cigarettes. That is, the machine dispenses the selected video cassette by dropping the cassette from a column of like titles to a platform located below the column where the person making the selection can remove the cassette from the machine. After completing the use of the cassette, the customer returns the cassette by depositing it in a storage bin associated with the machine. In order to position the cassette for subsequent dispensing, an operator must physically open the machine and place the cassettes contained in the storage bin in a respective one of the columns associated with the various movie titles.\nThe above described cassette vending machine has several disadvantages. As noted above, the design of the machine permits it to dispense only a limited number of titles. Second, the cassettes are subject to being damaged in the dispensing process when they are dropped to the platform and again when they are returned to the machine and dropped into the storage bin. Further, the machine requires constant oversight by an operator in order to store the returned cassettes in the appropriate position so that they can be later dispensed."}
{"text": "Some audiobooks and other items of audio content are recorded with only one voice actor. Multiple versions of a particular item of audio content may exist, each vocalized by only one voice actor. The one-voice approach may leave listeners unsatisfied, however. For example, consumers of an audiobook may want the female characters voiced by female voice actors, and male characters voiced by male voice actors. In another example, listeners of an audiobook may want a different voice actor for each character in the audiobook. Neither outcome is possible with only one voice actor.\nRecording a new audiobook with multiple voice actors is one option for customizing an item of audio content, but this approach has its limits. It may be impractical or expensive to coordinate multiple voice actors to generate a new recording. Additionally, listeners may want to assign voice actors to characters in different ways, which could necessitate a new recording for each listener. Moreover, listeners may have no interest in an audiobook featuring a disliked voice actor, even if multiple voice actors are already present in the audiobook. These problems and others are not merely limited to audiobooks, but may apply generally to any digital audio content, such as audio content included in songs, videos, movies, television programs, computer and video games, multi-media content, and the like."}
{"text": "Stone such as marble and granite are used for many purposes including counter tops, door sills, decorative inlays and the like. For many applications, the natural colors and random patterns are desired for aesthetic reasons.\nIn some applications, it is desired to impregnate the natural stone with a design or image. Such applications include decorative replacements for stained glass and photographic images in monuments, etc. Prior attempts at such have produced limited results with superficial images that wear with time.\nSome prior art includes methods of printing on stone. For example, U.S. Pat. No. 5,916,662 to Schmidt shows how to print on a coating on the stone. Unfortunately, the coating covers the stone and detracts from the aesthetic appeal of the stone and the coating can separate from the stone.\nU.S. Pat. No. 6,569,277 to Gibbs shows how to transfer an image onto the surface of a material including a leaf and stone. Unfortunately, placing the image on the surface results in an image that is easily scratched.\nU.S. Pat. No. 6,686,315 to Creed has a method of making a building material that simulates the look of marble or granite that may include lettering, etc. This method uses a coated substrate as in U.S. Pat. No. 5,916,662 and, therefore, does not present natural stone to the viewer.\nU.S. Pat. No. 7,108,890 to Horne, et al, also requires a coating or matrix to be applied to the stone before introducing the image and, therefore, does not present natural stone to the viewer.\nWhat is needed is a method of impregnating a stone material with an image that will augment the natural beauty of the stone with an indelible image."}
{"text": "1. Field of the Invention\nThis invention generally relates to a lens barrel and more particularly, to a lens barrel that has a light shielding (namely, intercepting or blocking) portion for preventing light rays from being incident through an opening between barrels which are adapted to move relatively in a direction of an optical axis and/or in a direction of rotation thereof.\n2. Related Art Statement\nA camera provided with zooming and collapsing functions has a lens barrel consisting of a plurality of barrels such as a lens cylinder, a linearly (or rectilinearly) movable (or moving) cylinder, a rotary (or rotating) cylinder and a stationary (or fixed) cylinder. Such a type of camera realizes (or implements) the aforementioned zooming and collapsing functions by relatively displacing the aforesaid barrels or displacing the aforementioned barrels relative to an outer cover for a camera main body in a direction of an optical axis and/or in a direction of rotation thereof.\nIn the aforementioned lens barrel, small gaps are provided between the barrels facing each other and between the barrel and the outer cover for the camera main body, respectively, so as to reduce a drive (or driving) load at the time of relatively displacing the barrels. If adequate measures are not taken to block light at the gaps, there is caused a defect in that unnecessary external (or extraneous) light enters directly the main unit of a camera or indirectly therein (after undergoing secondary reflection or scattering) and then reaches a film surface and thus a film is exposed to the light.\nTherefore, there has been made various propositions for preventing external light from entering the camera main body through the aforementioned gaps. For example, in Japanese Unexamined Utility Model Publication No. 5-61711 Official Gazette, there has been proposed the following light shielding device for use in a camera.\nNamely, as shown in FIG. 10, the aforesaid light shielding device has a light shielding member 103 formed by sticking (or putting) a ring-like synthetic resin sheet member 103a, whose inside (or inner) diameter is less than the outside (or outer) diameter of a movable cylinder 101 for holding a (picture) taking lens, onto a ring-like light shielding (or intercepting) sheet base plate 103b. Further, the aforesaid light shielding member 103 is fixed to a front cover 102 of the camera main body. Then, the aforementioned sheet member 103a is pressure-connected (or pressure-contacted) to and is brought into surface contact with the aforementioned movable cylinder 101. Thus, external light is prevented from entering the camera main body through the gap between the movable cylinder 101 and the front cover 102 thereof.\nThere have been employed light shielding devices, most of which have the configuration as above described.\nMeanwhile, the light shielding device for use in a camera, which is proposed in the aforesaid Japanese Unexamined Utility Model Publication No. 5-61711 Official Gazette, is a mechanism by which the shielding of light is achieved by bringing only a single elastic (or resilient) light shielding means (corresponding to the sheet member 103a in the aforementioned example of the device) into surface contact with the surface of the inner wall of the movable cylinder. However, the light shielding device of such a configuration has encountered problems that the surface contact is unstable in the case where the aforesaid movable cylinder is eccentric with respect to an optical axis and that thus the complete light shielding performance cannot be obtained.\nOn the other hand, in recent years, there have been increasing demands for miniaturization of a camera. This exerts an influence on barrels. Further, when miniaturizing the barrels, there arises a necessity for miniaturizing a light shielding device incorporated in a camera.\nHowever, when miniaturizing the aforementioned light shielding device, there is a fear that the following drawbacks, especially, deterioration in light-shielding performance may be caused.\nTo begin with, drawbacks caused in the case of reducing the size in a direction (namely, in a radial direction), which is perpendicular to an optical axis, of the aforementioned elastic light shielding means will be described hereinbelow.\nFirst, in the case where it is intended that the miniaturization is achieved by reducing the outside diameter of a supporting member (corresponding to the light shielding sheet base plate 103b) for the elastic light shielding means (corresponding to the sheet member 103a), the mounting strength of the aforesaid sheet member decreases. Thus, the contact of the aforesaid sheet member with the movable cylinder 101 becomes unstable. Consequently, there is caused a drawback in that the complete light shielding performance cannot be obtained.\nNamely, when reducing the outside diameter of the aforesaid supporting member, the stiffness (or rigidity) of the aforesaid supporting member itself decreases. The relative displacement between the barrels causes deformation of the supporting member, which receives a force exerted in a direction of an optical axis, into a cone-like (or conical-spring-like) form. Thus, an amount of lap in the sheet member decreases. Consequently, there is caused a drawback in that unnecessary light enters the camera main body.\nSecond, in the case where it is intended that the miniaturization is attained by reducing the inside diameter of the aforementioned supporting member without changing the thickness in the radial direction thereof and by further cutting a lap portion of the aforementioned supporting member, the amount of the remaining lap may be small according to the quantity of the cut lap. In this case, owing to the eccentricity of the barrel (namely, the movable cylinder), there is caused a drawback in that external light is liable to leak into the camera main body.\nMoreover, the curvature of a bending (or bent) portion of the sheet member decreases. Thus, there is caused another drawback in that a load imposed on the barrel increases. Namely, even if the sheet member is formed in such a manner as to be thin to some extent, a rather large load is produced between the barrel and the sheet member, with the result that the drive load between the barrels increases.\nNext, drawbacks caused in the case of reducing the size in the direction of the optical axis (or the optical axis direction) of the aforementioned elastic light shielding means will be described hereunder.\nFirst, in the case where it is intended that an amount of lap in a portion in the optical axis direction of the aforementioned sheet member 103a, at which the sheet member 103 is in contact with the movable cylinder 101, is reduced, there is caused a drawback in that external light is liable to leak into the camera main body, similarly as in the aforementioned case. Further, the curvature of a bending portion of the sheet member decreases. Thus, similarly as in the aforementioned case, there is caused another drawback in that a load imposed on the barrel increases.\nSecond, in the case where it is intended that the miniaturization is achieved by reducing the sheet thickness of the sheet member, the light blocking quality of the aforesaid sheet member itself is degraded. Namely, in the case where the sheet member is formed in such a manner that the sheet thickness thereof is small, external light is transmitted by the sheet member when the external light impinges directly on the sheet member. In addition, when light rays are transmitted by the sheet member, the light rays are scattered. Thus, there is caused a drawback in that a wide range of a film surface is exposed to the light rays having been incident in a camera.\nThird, in the case where it is intended that the miniaturization is attained by reducing the thickness of the supporting member for supporting the sheet member, the stiffness of the sheet member is degraded. Thus, the contact of the sheet member with the movable cylinder becomes unstable. Consequently, similarly as in the aforementioned case, there is caused another drawback in that the complete light shielding performance cannot be obtained."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a ratchet wrench and, more particularly, to a thin or low-profile ratchet wrench for socket driving or removal having a socket reversing lever projecting from a terminal end of the ratchet wrench handle for adjusting the direction of socket driving or removal.\n2. Description of the Related Art\nWhen the automobile industry shifted production emphasis to economy size automobiles, the size of the engine and the concomitant space were significantly reduced. The reduction of size and space has carried over to mid-size and large vehicles. Although the reduction in engine size and space lowered production costs, and ultimately consumer prices, the ability of the do-it-yourself mechanic to perform maintenance on an auto was permanently hindered. The bolts, screws and nuts used on an engine are placed in spaces that are quite difficult to access with ordinary hand tools. Similarly, certified mechanics were required to purchase new tools and spend more time in disassembly and reassembly because of the compact space, increasing the labor costs to consumers. Thus, a need for an inexpensive hand tool appropriately sized to accommodate use in tightly spaced areas is apparent.\nA search of the prior art did not disclose any patents that read directly on the claims of the instant invention; however, the following references were considered related:\nU.S. Pat. No. 6,536,310, issued in the name of Goldfein, discloses a low profile wrench;\nU.S. Pat. No. 6,253,647, issued in the name of Eggert et al., discloses a reversible ratchet with remote reversing operating mechanism;\nU.S. Pat. No. 6,234,049, issued in the name of Ju et al., discloses an improved ratchet tool with a swing head that is rotated within a small angle by an eccentric member connected to a rod in the shank;\nU.S. Pat. No. 6,148,694, issued in the name of Spirer, discloses a hand wrench with torque augmenting means;\nU.S. Pat. No. 5,857,390, issued in the name of Whiteford, discloses a reversible ratchet wrench having an axial opening to allow a bolt to pass through and extend beyond;\nU.S. Pat. No. 5,584,220, issued in the name of Darrah et al., discloses an angle attachment tool having a drive element at a remote end to accommodate a power tool;\nU.S. Pat. No. 4,907,476, issued in the name of Singleton, discloses a speed socket wrench with a T-shaped turning handle;\nU.S. Pat. No. 4,311,072, issued in the name of Hudgins, discloses a speed handle ratchet wrench;\nU.S. Pat. No. 2,808,749, issued in the name of Lampke, discloses a gear operated power wrench; and\nU.S. Pat. No. 2,549,515, issued in the name of Orey et al., discloses a wrench having a pivoted handle and removable socket.\nConsequently, there is a need for a thin or low-profile ratchet wrench with a socket reversing lever projecting from an end of the ratchet wrench handle."}
{"text": "1. Field of the Invention\nThis invention relates in general to the field of closing devices useful in preventing water from passing to the chain locker via the hawse pipe of a marine vessel. Specifically the invention relates to a sealing device, which when installed in the hawse or chain pipe leading to the chain locker of a drilling vessel, will seal about the chain to prevent flooding of the chain locker during extreme list or storm weather conditions.\n2. Description of the Prior Art\nThe prior art has provided both horizontal and vertical arrangements for the hawse pipe leading to the chain locker of a semi-submersible rig. Where the winch is located in the chain locker below the main deck, the top of the platform is left relatively clear and the overall center of gravity of the rig is advantageously lowered. If the hawse pipe is horizontally oriented, the hawse pipe opening may be low in relation to the water but, the rig may be susceptible to flooding during extreme list conditions or heavy weather.\nIn the vertical arrangement of the hawse pipe, a deck sheave and an underwater fairlead direct the anchor chain from an anchor to the vertical hawser pipe leading to the chain locker. Even with the vertically oriented hawser pipe, the chain locker is a potential rig destabilizing problem if it were to fill with water during extreme list conditions or heavy weather.\nThe prior art has not been able to solve the problem of providing a reliable sealing device to seal about the hawser or anchor chain and prevent seawater from flowing to the chain locker from the hawse pipe.\nPrior art attempts to seal between the hawse pipe and the hawser chain have been unsuccessful due to the inability to position the chain relative to a ram-type closing device and due to the inability of other types of closing devices to conform to the highly irregular profile of a chain."}
{"text": "In an enterprise, various pieces of content may be scattered around the enterprise. For example, documents may be stored in various repositories, may be shared between individuals, and may be interacted with by various individuals via a variety of workloads. Some content items may be interesting to or relevant to a particular user; however, he/she may not be aware of why the content may be interesting to him/her or that he/she may have a personal connection to it.\nA content item may have been viewed, edited, commented on, etc. by people with whom the user has a relationship. A traditional document browsing or search experience may provide regular metadata such as a content item's title, location, file size, summary, etc., but may not provide insights as to why the content item may matter to him/her. It is with respect to these and other considerations that the present invention has been made."}
{"text": "The present invention relates to methods and computer readable media for preserving unique critical information during data imaging.\nProviding backup data to a hard drive of a computing device is of vital importance in today\"\"s information age. Individuals have become so dependent upon using their computing device, that they have become non functional when the device is inoperable. When a computing device starts up, or is booted up, it uses a hard drive memory associated with the computing device to provide all the necessary data files and data application programs to the user of the device. Some of these files may reside locally on the computing device, while others are provided remotely to the computing device if the computing device is networked with other devices providing files and applications.\nA standard process, referred to as hard drive imaging, allows the user to take a snapshot of a hard drive which includes all data included therein and store these data in an image archive. This image archive can be stored off onto additional computer readable media, such as and by way of example only, CDs, diskettes, additional hard drives, and other magnetic or optical storage devices. The image archive can later be used to restore the original computer (and/or any other reasonably compatible device) to the state in which it existed at the time the image archive was created.\nDuring the imaging process all information is captured exactly and precisely. Should a failure occur with a user\"\"s hard drive on the original computer or on a compatible device, the image archive can be used to place all data contained on the original device at the time the image archive onto the current device. This process effectively places the device into the exact state that existed on the original device at the time the image archive was created. Unfortunately, the state duplication is too exact for some applications and network functions.\nSome applications and network functions require that certain information be unique among all devices attached to the same network as the current device. For example, TCP/IP networking requires that all addresses on a network be unique. Other networking applications require that each device have a unique name. Yet other applications may require that certain data be the same on a certain administrator-defined group of devices, yet different from all other defined groups.\nThe results for failing to set these values correctly after an image archive has been restored to the device vary from application to application. Some cause the network components of the device to function incorrectly, or to shut down, resulting in the complete loss of access to the network. Others may results in the user receiving different security access privileges than intended. Results are almost impossible to predict and vary greatly from application to application, but all result in the need for a manual reconfiguration of the device after the image archive has been restored.\nQuite often, a device needs to be restored using a different image archive or a similar image having modified identifying information, because some catastrophic event has occurred to make it inoperable. Prior to that event, however, this device would have had valid values for all data requiring any degree of uniqueness. Unfortunately, these values are irretrievably lost during the hard drive imaging process using the different or similar image archive.\nAccordingly, an object of the invention is to provide methods and computer readable media to preserve unique critical information during the imaging process and to restore it after the image process has completed. This information is referred to collectively as image-safe data and includes, but is not limited to static IP addresses, workgroup/role identifying information, file permissions, and the like. Using the image archive and the image-safe data a computing device may be effectively restored to a fully functional state, should its main memory (e.g. hard drive) crash, or otherwise become corrupted.\nThis image-safe data may be stored directly on non-partitioned and unused areas of the hard drive of a user\"\"s computing device, with the typical image archive stored on removable or remote computer readable media. In this way, a user may restore a computer device to a fully functional state using the image safe data recorded on the unusable portions of the hard drive (e.g., non-partitioned areas) and the image archive stored on the removable computer readable media.\nShould a hard drive failure occur, even a user not currently connected to the network could be restored to a fully functional state by using the removable media as a boot device. This would then permit the user to reformat the computing device\"\"s hard drive and restore all information properly, including the image safe data. With the image safe data, the user may then successfully connect to a network and identify itself properly or acquire the appropriate programs based on the user\"\"s defined work groups, roles, file permissions, and the like.\nAdditional objectives, advantages and novel features of the invention will be set forth in the description that follows and, in part, will become apparent to those skilled in the art upon examining or practicing the invention. The objects and advantages of the invention may be realized and obtained by means of the instrumentalities and combinations particularly pointed out in the appended claims. To achieve the foregoing and other objects and in accordance with the purpose of the present invention, methods and computer readable media are provided for preserving critical information during an imaging process.\nA method of augmenting the imaging process of a computer device\"\"s memory is provided, comprising the executable instructions of acquiring information necessary to properly operate one or more network applications or uniquely identify the computing device to a network of additional computing devices. Further, from time to time the information is acquired in one or more storage locations of the computing device, wherein the locations are not partitioned or otherwise used by the computing device.\nIn another aspect of the present invention, a computer readable medium is provided, having one or more usable locations wherein the usable locations are operable to be partitioned and used without any significant restrictions. Moreover, one or more reserved locations are used to house information necessary to properly operate one or more network applications or uniquely identify a specific computing device. The reserved locations may not be partitioned and are only accessed by authorized applications.\nIn yet another aspect of the present invention, a method of restoring a memory associated with a computing device connected to a network of computing devices is provided, having executable instructions wherein from time to time data associated with a memory of a computing device is acquired. Furthermore, from time to time one or more non-partitioned memory locations identifying information which uniquely identifies the computing device is acquired. Upon a request, the data and the identifying information are restored."}
{"text": "1. Field of the Invention\nThe present invention relates to a printing apparatus which performs printing based on print data sent from a host computer.\n2. Description of the Related Art\nThere is conventionally known a printing apparatus (label writer) in which printing on a tape is performed based on print data or control data which is sent from a host computer and the tape is cut to have a predetermined length, thus creating a label (see Published Unexamined Japanese Patent Application No. 2002-036397, FIG. 1 and related description). This type of printing apparatus does not have a user interface on a printing apparatus case thereof for creating and editing data. Therefore, when conducting shipping inspections of products or performing demonstration printing in sales activities, printing has to be carried out by sending print data to the printing apparatus from the host computer after a printing apparatus driver and a print-only application are installed into the host computer, and thereafter the host computer and the printing apparatus are connected using a cable or the like.\nDue to this, when performing demonstration printing, it is required to prepare a display, a keyboard, a cable and the like, in addition to the host computer in which the printing apparatus driver and the like are installed. Therefore, preparation thereof requires time and effort. In addition, in sales activities, it has been very inconvenient in that the host computer had to be carried around or a host computer with an operation system (OS) which corresponds to the printing apparatus had to be prepared by the customer."}
{"text": "This invention generally relates to electrical component connectors and contacts. More particularly, the present invention relates to an inventive three point contact design that may be used to connect varied electrical components to circuit boards such that the components may be assembled or installed on the circuit board and removed from the board without the need for time consuming soldering and desoldering of the component leads. The inventive contacts use flexible integral spring elements and appropriate contact element shaping to ensure a solid electrical connection between the printed circuit board and the electrical component leads. With the use of the inventive spring contacts, components may be easily installed on a board and tested, and, if necessary, may later be removed or replaced without the need for time consuming soldering and desoldering of the component. The inventive contacts are designed to ensure sound electrical connection with component leads of different cross sections, shapes and sizes.\nAs the size of electrical components used in electronic circuitry continues to decrease in size, the density of circuitry placed on printed circuit boards increases. With these two trends, a decrease in overall size and an increase in circuitry density, the value of space or real estate on circuit boards has similarly significantly increased. In addition to the importance and value of space on the board, equally important is need to ensure sound electrical connections between the component leads, which are decreasing in size and structure, and the circuit board.\nAnother trend in the design and manufacture of circuit boards that kept pace with the complexity and number of circuits and components included on boards, is the need to test, and potentially replace, the fabricated circuits and installed components as they are fabricated or installed on the board and tested. The need to disassemble or replace components on circuit boards has become an important issue. For example, it is inefficient and costly to install a component on a board, to only find out, after the board assembly is complete, that the component is not functionally performing. The time and effort to disassemble, remove and replace the component is expensive and fraught with the potential for creating or causing additional problems in the circuit.\nOne aspect of this problem is exemplified with electrical components that are installed on a board by soldering of the component leads to the board contacts. The time and effort to install a component, solder the leads to the contacts, test the component and, if the component is not functionally operable, desolder the leads and replace the component is very inefficient and expensive. As noted, in each of the desoldering, disassembly, reassembly and resoldering steps, there is the potential for creating additional problems in the board fabrication.\nThe current designs for installing electrical components to a board are similar to integrated circuit sockets, where the electrical component leads are bent 90 degrees so that the lead may be inserted into the socket and soldered in place. Again, the soldering of the leads to ensure a sound electrical connection does not allow for quick or easy disassembly and removal of the component.\nAccordingly, there remains a need for a device or contact that allows an electrical component to be easily installed onto a circuit board without the need for solder connections, that ensures sound electrical connections and paths between the component leads and the circuit board, that allows the component to be tested through the board electrical connections, and also allows the component to, if necessary, be easily removed from the board and replaced with the need for desoldering the component leads.\nSuch a device and contact would greatly increase the efficient assembly and fabrication of circuit boards necessary for mass production of electrical component packages used in various electronic devices.\nIn view of the shortcomings of the prior art, it is an object of the present invention to provide an electrical contact that allows an electrical component to be easily and quickly installed on a circuit board and has sound electrical connections between the component and the board without the need for soldering of the component leads to the board contacts. It is a further object of the present invention that the electrical contact allows an electrical component to be easily and quickly removed from a circuit board without the need for desoldering of the component leads or without damaging the board contacts.\nTo achieve this and other objects, and in view of its purposes, the present invention provides an electrical spring contact for use with varied electrical components, the spring contact comprising a body section, where the body section is attachable at one end to a circuit board; two arms attached to opposite sides of the body section, the two arms defining a separation between the arms; the two arms each having an elbow section extending each arm towards the opposite arm; and a center section attached to the body section between the two arms whereby an electrical component lead can be held by and between, and be in electrical contact with the center section and the two arms.\nIt is a further object of the present invention to provide an electrical spring contact for use with varied electrical components, the spring contact comprising a body section having a first and second end and two opposing sides, the body section being attachable at the first end to a circuit board; two arms attached to the opposing sides of the body section and extending away from the second end of the body section, the two arms defining a separation between the arms, the arms each having an elbow section such that the separation between the two arms is reduced downstream of the elbow sections; and a center section attached to the body section between the two arms whereby an electrical component lead can be held by and between, and be in electrical contact with the center section and the two arms.\nIt is a further object of the present invention to provide an electrical spring contact for use with varied electrical components, where the electrical spring contact can hold an electrical component lead having a diameter approximately in the range of 0.010 inches to 0.025 inches.\nAnother aspect of the present invention is an electrical component assembly, having at least one electrical component, the electrical component having at least one electrical lead, and the connector assembly comprising a housing into which the electrical component fits and is held; and at least one electrical spring contact held within the housing, said electrical spring contact comprising, a body section, the body section being attachable at one end to a circuit board; two arms attached to opposite sides of the body section, the two arms defining a separation therebetween, the arms each having an elbow section extending each arm towards the opposite arm; and a center section attached to the body section between the two arms whereby the at least one electrical component lead can be held by and between, and be in electrical contact with the center section and the two arms.\nIt is a further object of the present invention to provide an electrical contact assembly, having at least one electrical component, the electrical component having at least one electrical lead, the contact assembly comprising a first housing into which the electrical component fits and is held; a second housing; at least one electrical spring contact seated within the second housing, the electrical spring contact comprising, a body section, the body section being attachable at one end to a circuit board; two arms attached to opposite sides of the body section, the two arms defining a separation therebetween, the arms each having an elbow section extending each arm towards the opposite arm; and a center section attached to the body section between the two arms whereby the at least one electrical component lead can be held by and between, and be in electrical contact with the center section and the two arms.\nThese and other aspects of the present invention are set forth below with reference to the drawings and the detailed description of certain preferred embodiments. It is to be understood that both the foregoing general description and the following detailed description are exemplary, and are not intended to be or should be considered restrictive of the invention."}
{"text": "The present invention relates generally to communication systems and, in particular, to a communication system that employs a method and apparatus for allowing a user of a display-based terminal to communicate with communication units in the communication system.\nDispatch communication systems are known to include a wireless infrastructure and a plurality of communication units, such as two-way radios. Some dispatch systems, such as those used for public safety, also include a dispatch console and a so-called xe2x80x9ccomputer aided dispatch (CAD) systemxe2x80x9d that includes a display-based terminal to control communications between the communication units. The CAD terminal typically displays categorized tables of information to the terminal user (typically referred to as a xe2x80x9cdispatcherxe2x80x9d or xe2x80x9cdispatch operatorxe2x80x9d). For example, the CAD terminal may display queues of currently pending incidents and/or a list of communication units that are currently available. In addition, some CAD terminals include an integrated mapping program that enables the CAD terminal to display locations of communication units on a map that represents a geographic area supported by the dispatch system. The locations of the communication units are typically provided to the CAD system on a periodic basis by an automatic vehicle location (AVL) system that is coupled to the CAD system via a dedicated communication link.\nIn addition to receiving communication unit location information, the CAD system may also receive incident information from a 911 system that is coupled to the CAD system. In such a case, the map displays the origination point of a 911 telephone call to the dispatcher as an icon on the map. By viewing communication unit location, communication unit status, and incident location on the map, the dispatcher can quickly determine which communication unit users (e.g., policemen, firemen, paramedics, and so forth) would be in the best situation to respond to the incident. As an incident is attended to by users of the communication units, the status of the communication units associated with such users is updated either manually by the dispatcher or automatically by the CAD system responsive to messaging from the wireless infrastructure. In the latter case, the wireless infrastructure receives status updates over a wireless communication channel from the communication units that are participating in the handling of the incident. Changes in communication unit status are typically indicated to the dispatcher by some type of visual change, such as a color change or icon update, to the participating unit\"\"s representation on the map.\nA map display, with icons representing the incidents and the communication units with their current status, gives a dispatcher a powerful tool for quickly evaluating an emergency situation. By using standard graphical user interface (GUI) cursor interaction (e.g., point and click or drag and drop), the dispatcher can manipulate the screen icons to assign units to incidents. In some systems, the action of assigning communication units to incidents on the map automatically sends data to the assigned communication units (e.g., if the units are equipped with mobile data terminals) to enable the communication unit users to immediately obtain critical information related to the incident, such as the incident location and other details.\nIn many situations, the assignment of a communication unit to an incident is followed by a voice conversation between the dispatcher and the user of the communication unit. Such a voice conversation enables the dispatcher and the communication unit user to discuss additional details about the incident which may not have been automatically relayed via data messaging. Voice communication during the handling of an incident is critical because it provides the communication unit users with real-time updates of the situation and the parties involved as the incident progresses.\nAlthough the aforementioned voice communication between the dispatcher and the assigned communication unit users is critical, such voice communication cannot be initiated from the map environment with prior art CAD systems or terminals. Rather, in existing systems, the dispatcher must access a separate dispatch console program (which may be displayed on the CAD terminal in a different environment from the map) to control the voice communication with the communication unit users, thereby losing the benefit of viewing the overall picture of the incident as provided by the map.\nOne conventional dispatch system is currently being developed for a municipal transit authority, but as yet has not been implemented, in an attempt to integrate voice control into the CAD terminal map display. In this proposed system, which will include a first map displaying a subway system and a second, separate map displaying train locations, the dispatcher will be able to select a section of the subway from the first map and automatically initiate a voice call to all trains in the selected area of the subway only. As noted above, the first map display does not, however, include the locations of the trains. Therefore, the dispatcher cannot select particular trains individually or as a group to communicate with. Moreover, the proposed system does not facilitate trunked communications and does not take into account various communication parameters, such as private line, digital private line, encryption, and others. Therefore, although this proposed system may allow initiation of voice calls to a geographic area depicted on the first map, the system does not insure that any communication will actually occur because the dispatcher will not know from the first map display whether or not trains are in the selected area of the subway and the dispatcher cannot select the trains themselves, either individually or in a group, to communicate with from either map display. Consequently, although the proposed dispatch system does suggest a type of map-based voice dispatch, it has many practical shortcomings.\nTherefore, a need exists for a method and apparatus that allows a user of a display-based terminal, such as a CAD terminal, having an integrated mapping program to communicate directly with communication units from a single map environment."}
{"text": "The use of touch-sensitive surfaces as input devices for computers and other electronic computing devices has increased significantly in recent years. Example touch-sensitive surfaces include touchpads and touch-screen displays. Such surfaces are widely used to manipulate user interface objects on a display.\nExample manipulations include adjusting the position and/or size of one or more user interface objects or activating buttons or opening files/applications represented by user interface objects, as well as associating metadata with one or more user interface objects or otherwise manipulating user interfaces. Example user interface objects include digital images, video, text, icons, control elements such as buttons and other graphics. A user will, in some circumstances, need to perform such manipulations on user interface objects in a file management program (e.g., Finder from Apple Inc. of Cupertino, Calif.), an image management application (e.g., Aperture, iPhoto, Photos from Apple Inc. of Cupertino, Calif.), a digital content (e.g., videos and music) management application (e.g., iTunes from Apple Inc. of Cupertino, Calif.), a drawing application, a presentation application (e.g., Keynote from Apple Inc. of Cupertino, Calif.), a word processing application (e.g., Pages from Apple Inc. of Cupertino, Calif.), a website creation application (e.g., iWeb from Apple Inc. of Cupertino, Calif.), a disk authoring application (e.g., iDVD from Apple Inc. of Cupertino, Calif.), or a spreadsheet application (e.g., Numbers from Apple Inc. of Cupertino, Calif.).\nBut methods for performing these manipulations are cumbersome and inefficient. For example, using a sequence of mouse based inputs to select one or more user interface objects and perform one or more actions on the selected user interface objects is tedious and creates a significant cognitive burden on a user. In addition, these methods take longer than necessary, thereby wasting energy. This latter consideration is particularly important in battery-operated devices."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic resonance imaging apparatus having a user interface capable of easily determining a tomographic position and a method of controlling the same.\n2. Description of the Related Art\nA magnetic resonance imaging apparatus creates images of chemical and physical microscopic information of materials or observes a chemical shift spectrum by using a phenomenon in which a group of nucleuses having an inherent magnetic moment resonantly absorbs energies of a high frequency magnetic field rotating at a specific frequency when they are located in a uniform static magnetic field.\nIn a tomography using the magnetic resonance imaging apparatus, it is generally necessary to set various tomographic parameters (tomography sequence) such as, for example, TR, TE, a slice thickness, a size of FOV or the like. Since each tomographic parameter has an influence on the contrast of an image correlating with the other parameters, some technical knowledge or technologies and experiences are required to obtain images suitable for a diagnosis. In general, an operator performs another modality. For this reason, a period when the magnetic resonance imaging apparatus is not operated is generated, so that it is difficult for the operator to get used to the operation of the magnetic resonance imaging apparatus. As a result, there are cases that the operation time grows longer and the throughput of the work degrades.\nIn order to solve the above-mentioned problems, there are technologies called a simple tomography or PAS (programmable anatomical system). The simple tomography is a tomography method for beginners in which the beginners can relatively easily tomograph a MR image. In addition, the PAS is a function that can automatically preset tomographic parameters for each object part before scanning. By using these technologies, even engineers not having sufficient experiences can set suitable tomographic parameters to tomograph a suitable diagnosis image.\nHowever, in the magnetic resonance imaging apparatuses in the related art, setting a position and an angle of a ROI (region of interest) has the following problems. Specifically, in the related art, most of the tomographic parameters can automatically be set by the PAS technology. However, the position and angle of the ROI cannot be automatically set. As a result, the operator subjectively determines and sets the position and angle of the ROI through the manual operation with a sufficient attention. At this time, a user interface setting the position and angle of the ROI from the positional relationships between an outline and an internal organ in a subject is used. However, there is a possibility that the unaccustomed user sets the ROI at the erroneous position and angle and tomographs it.\nFurthermore, since the magnetic resonance imaging apparatus can tomograph all parts of the human body, a considerable training time is needed to master the flawless tomography technique. Therefore, in order to solve these problems through mastering the specialized knowledge or technology, a burden is large from a viewpoint of time, cost and management. In addition, from a medical point of view requiring an early diagnosis and therapy, the magnetic resonance imaging apparatus is not necessarily preferable."}
{"text": "1. Field of the Invention\nThe present invention relates to a non-volatile memory device and, more particularly, to methods of erasing a non-volatile memory device having discrete charge trap sites.\n2. Description of the Related Art\nA non-volatile memory device is configured to retain data previously stored in its memory cells, even when the power supply is turned off, unlike a volatile memory device which loses stored data when power is removed. The non-volatile memory device includes a charge storage layer, which is interposed between a semiconductor substrate and a gate. Non-volatile memory devices may be classified as a floating gate type memory device and an MNOS type memory device, depending on the structure of the charge storage layer.\nThe floating gate type non-volatile memory device structure includes a tunnel dielectric layer, a floating gate, an inter-gate insulating layer, and a control gate, which are sequentially stacked on a semiconductor substrate. The floating gate, in which the charge is stored, is formed of a conductive layer.\nThe MNOS type non-volatile memory device structure includes a stacked metal/nitride/oxide/semiconductor (MNOS) or metal/oxide/nitride/oxide/semiconductor (MONOS). In this configuration, a dielectric layer such as nitride, is interposed between a semiconductor substrate and a gate, to serve as a charge storage layer. The MNOS type non-volatile memory device stores data by using trap sites within the dielectric layer, and trap sites in an interface of the dielectric layer, for example, the interface between adjacent dielectric layers, and trap sites in an interface between the dielectric layer and the semiconductor substrate.\nChan, et. al. introduced SONOS (silicon oxide nitride oxide silicon)-type memory devices in “A True Single-Transistor Oxide-Nitride-Oxide EEPROM Device, IEEE Electron Device Letters, Vol. 8, No. 3, pp. 93–95, 1987). The conventional SONOS type non-volatile memory device includes a dielectric layer and a gate, which are formed on a silicon substrate, and source/drain formed in the semiconductor substrate adjacent to both sidewalls of the gate. The dielectric layer includes a tunnel dielectric layer, a charge storage dielectric layer, and a blocking dielectric layer. The tunnel dielectric layer, the charge storage dielectric layer and the blocking dielectric layer may be formed of a silicon oxide layer, a silicon nitride layer and a silicon oxide layer respectively. The gate may be formed of a silicon layer. That is, the SONOS structure is formed of the silicon oxide layer, the silicon nitride layer, the silicon oxide layer and silicon layer, which are sequentially stacked on the silicon substrate.\nWhen the SONOS type non-volatile memory device becomes programmed, and electrons are injected into the charge storage dielectric layer, the device threshold voltage is increased. Thus, when a voltage, being lower than the increased threshold voltage, is applied on the gate, electric current does not flow through the channel. When a voltage higher than the threshold voltage is applied, current flows through the channel. Using this property, stored data can be read.\nFurthermore, the erase operation of the SONOS type non-volatile memory device is performed by applying a negative voltage to the gate and by applying a positive voltage to the source, while floating the drain and the semiconductor substrate. As such, the electrons stored within the charge storage dielectric layer are removed. In the case where the charge storage layer is a floating gate comprising a conductive layer, the electrons stored within the charge storage layer can be freely moved. On the contrary, in the case where the charge storage layer is formed of a dielectric layer, such as in the case of the SONOS type non-volatile memory device, the electrons stored inside the trap sites cannot be freely moved. Thus, the conventional erasing method described above fails to provide a complete erasure in such a SONOS type non-volatile memory device.\nFIG. 1 is a graph showing the change of threshold voltage as a function of the number of program(P)/erase(E) cycles in the typical SONOS type non-volatile memory device. As shown in FIG. 1, while program threshold voltages remain somewhat constant as a function of the number of program/erase operations, the erase threshold voltages are gradually increased. Such an increase of the erase threshold voltages results from the fact that the electrons trapped within the charge storage layer are not completely removed as a result of the erase operation, as described above. If the increase in erase threshold voltage cannot be suppressed, after approximately 100,000 P/E cycles, a threshold voltage window between a programmed non-volatile memory device and a non-programmed non-volatile memory device is not obtained, thereby resulting in deterioration of device reliability."}
{"text": "Molecular interactions are critical for maintaining cellular functions. These interactions can be classified into 3 broad categories. The first is protein complex formation via covalent bonds such as the heavy chain and light chain of an immunoglobulin. The second is protein-protein association via non-covalent bonds exemplified by heterodimer formation of HER-2 and EGFR (epidermal growth factor receptor). The third is an association between 2 molecules involved in a cellular pathway such as a cytokine receptor and a caspase in an apoptosis pathway.\nAntigen binding proteins involved in the immune response are present in mammals as large polyclonal repertoires representing a broad diversity of binding specificities. This diversity is generated by rearrangement of gene sequences encoding variable regions of these binding proteins. Such variable region binding proteins include soluble and membrane-bound forms of the B cell receptor (also known as immunoglobulins or antibodies) and the membrane-bound T cell receptors (TCR). With respect to immunoglobulins, their affinity is enhanced subsequent to recognition of an antigen by a B cell antigen receptor, through a process termed affinity maturation which involves cycles of somatic hypermutation of these variable genes.\nKnown approaches for isolating antibodies with a desired binding specificity most often either involves generation of hybridoma cells from immunized hosts followed by screening for specific clones or involves the generation of combinatorial expression libraries in E. coli composed of immunoglobulin variable domains, which are subsequently screened using techniques such as phage display.\nThere are several limitations in the use of the hybridoma technology. The generation time for a specific hybridoma can be long (5-15 months). Functional screens are possible only after clone selection and culture. Furthermore, if a hybridoma is desired for producing human antibodies, typically for therapeutic purposes, then alternative strategies must be sought because of the absence of human myeloma lines suitable as fusion partners for human B lymphocytes. Heterohybridomas, i.e., fusion of human B cells with mouse myeloma cells, were attempted but they are extremely unstable and thus rarely lead to suitable cell lines for production purposes. Human B cells immortalized through infection with Epstein-Barr virus exhibit similar instability.\nUse of combinatorial libraries and phage display allows for generation of large repertoires of antibody clones with a potential diversity in excess of 1010. From this repertoire, selection for binding to a specific target can be performed thereby generating a sub-library. This sub-library can be used to generate either polyclonal or monoclonal antibodies. The variable region encoding sequences (for example heavy chain variable region and light chain variable region encoding sequences) which are to constitute the library can be amplified from lymphocytes, plasma cells, hybridomas or any other immunoglobulin expressing population of cells. Current technologies for generating combinatorial libraries involve separate isolation of the variable region encoding sequences. Thus, the original pairing of, for example, immunoglobulin heavy chain variable region and light chain variable region encoding sequences is lost. Said sequences are randomly paired and the original combinations of these variable sequences only occur by chance.\nThus, in order to isolate variable region encoding sequences responsible for a desired binding specificity, a considerable amount of screening is necessary. This is typically performed in combination with methods for enrichment of clones exhibiting a desired specificity, such as ribosome display or phage display. Even then, the diversity achieved might not be sufficiently large to isolate variable region encoding sequence pairs giving rise to binding proteins of similar high affinity as those found in the original cells. Further, the enrichment procedures normally used to screen combinatorial libraries introduce a strong bias, e.g. polypeptides of particular low toxicity in E. coli, efficient folding, slow off-rates, or other system dependent parameters, thus reducing the diversity of the library even further. In addition, clones derived from such combinatorial libraries will be more prone to producing binding proteins with cross reactivity against self-antigens because they, in contrast to original pairs (hereafter called cognate pairs), have never been subjected to in vivo negative selection against self-antigens as pairs, such as is the case for B- and T-lymphocyte receptors during particular stages of their development. Therefore, the identification of cognate pairs of variable region encoding sequences is a desirable approach. Moreover, the frequency of clones exhibiting a desired binding specificity is expected to be considerably higher within a library of cognate pairs than in a conventional combinatorial library, particularly if the starting cells are derived from a donor with high frequency of cells encoding specific binding pairs, e.g., immune or immunized donors.\nIn order to generate cognate pair libraries the linkage of the variable region encoding sequences derived from the same cell is required. At present, 3 different approaches that achieve cognate pairing of variable region encoding sequences have been described.\nIn-cell PCR is an approach where a population of cells is fixed and permeabilized, followed by in-cell linkage of heavy chain variable region and light chain variable region encoding sequences from immunoglobulins. This linkage can be performed either by overlap extension RT-PCR or by recombination. The amplification process, as described in these publications, is a three to four step process consisting of i) reverse transcription utilizing constant region primers generating immunoglobulin cDNA, ii) PCR amplification of the heavy and light chain variable region encoding sequences utilizing primer sets containing either overlap-extension design or recombination sites, iii) optional linkage by recombination, and iv) nested PCR of the products generating restriction sites for cloning. Since the cells are permeabilized there is a considerable risk that amplification products might leak out of the cells, thereby introducing scrambling of the heavy chain variable region and light chain variable region encoding sequences, resulting in the loss of cognate pairing. Therefore, the procedure includes washing steps after each reaction which makes the process laborious and reduces the efficiency of the reactions.\nMore generally, the in-cell PCR is extremely inefficient, resulting in low sensitivity. Accordingly, the in-cell PCR linkage technique has never found widespread usage, and the original studies have never been reliably repeated in a way that can verify that the linkage actually occurs within the cell. This, however, is absolutely crucial to avoid scrambling the heavy chain variable region and light chain variable region encoding sequences and thereby disrupting the cognate pairs.\nSingle-cell PCR is a different approach to achieve cognate pairing of heavy chain variable region and light chain variable region encoding. In these publications, a population of immunoglobulin expressing cells are distributed by diluting to a density of one cell per reaction, thereby eliminating scrambling of heavy chain variable region and light chain variable region encoding sequences during the cloning process. The process described is a three to four step procedure consisting of i) reverse transcription utilizing oligo-dT-, random hexamer- or constant region primers generating cDNA, ii) fractionating the cDNA product into several tubes and performing PCR amplification on the individual variable chain encoding sequences (in separate tubes) with primer sets containing restriction sites for cloning, iii) nested PCR of the products generating restriction sites for cloning (optional) and iv) linking the heavy chain variable region and light chain variable region encoding sequences from the separate tubes by cloning them into an appropriate vector, which is itself a multi-step process.\nIn humans, there are two types of light chains: lambda (λ) and kappa (κ). This means that with the cDNA generated from every single cell at least three separate PCR reactions must be performed followed by analysis and cloning of the appropriate bands into a single vector to achieve cognate pairing. Thus, the single-cell PCR approach as described requires a large number of reactions to generate a library of cognate pairs. Although, a cognate pair library does not need to be as large as a combinatorial library in order to obtain binding proteins representing a broad diversity of binding specificities it would still be a laborious task to generate a library of, for example, 104 to 105 clones by the described single cell PCR approach. Further, the large number of steps greatly increases the risk of contamination and human error.\nSymphogen's (Symphogen A/S, Lyngby, Denmark) approach is to use multiplex overlap-extension RT-PCR (reverse transcription-polymerase chain reaction) to identify the natural pairing of heavy chain and light chain antibodies by isolating single B cells in a well of a 96-well microtiter plate. However, the process is labor intensive and verification of the structure of individual antibodies is inefficient since each pair needs to be tested separately by expressing the corresponding single chain Fv in E. coli. \nPCR primer based DNA barcoding has been described in the literature. There have been studies utilizing oligonucleotide tags to label DNA molecules during sample preparation such that after sequencing by a massively parallel technology, one is able to digitally sort DNA sequences originating from different samples or a positive control. However, each tagging PCR primer must be individually synthesized. As such, when a large number of samples are processed, this practice will become prohibitively expensive. Consequently, the technique is only suitable for a limited number of samples."}
{"text": "(1). Field of the Invention\nThe present invention relates to cardiac monitoring and telemedicine systems and, more particularly, to cardiac monitoring systems which provide an analysis and display of one or more parameters relating to the condition of ischemic patients.\n(2). Description of the Related Art\nA number of new clog dissolving agents presented by the pharmaceutical industry during the past couple of years have given cardiologists the ability to immediately treat acute myocardial ischemia through chemical thrombolytic therapy. However, it is frequently difficult to properly control and adjust such therapy during the acute phase of a myocardial ischemia. Known methods are either expensive or have too large a delay (up to several hours) between the time of the myocardial ischemia time and the presentation of the results.\nSome cardiac monitoring systems and methods also utilize a known 12-lead electrocardiogram in which electrocardiogram (ECG) signals are displayed directly on a monitor in real-time. Such a 12-lead ECG arrangement has the disadvantages that a large number of electrodes must be placed on the patient in positions which cover mainly the frontal parts of the myocardium. A large storage capacity is also required in order to record all the ECG signals from the electrodes. However, many doctors are familiar with the format of the 12-lead ECG.\nFurther, the normal procedure in the case of serious illness is to transport the patient to a hospital for diagnosis and treatment of the illness. However, it has proved to be advantageous to arrange for the nursing or ambulance staff to carry out the diagnosis and to start the treatment already at the place of patient pick-up. Such an arrangement likewise makes it possible to establish at an early stage whether a particular specialist competence and special equipment or the like are required, whereafter the patient may be transported straight to the place where such competence, equipment etcetera are available.\nSuch early diagnosis and treatment would be considerably facilitated, were the nursing staff given a possibility to carry with them a portable telemedicine device adapted to register signals from ECG and vector electrocardiogram (VCG) units and similar sensing equipment. A portable unit of this kind, which may be docked and thus be connected to a stationary communications network (LAN), which allows connection thereto of external measurement equipment, and which comprises a display device for visualization of the measurement results, is disclosed in U.S. Pat. No. 5,375,604.\nHowever, this prior-art device is merely a passive unit and it is designed for reception and visualization only of signals from the measurement equipment. It cannot be used to establish active contact with and an exchange of information between the patient-attending staff and other individuals, a possibility which could be helpful in the diagnosis procedure as well as for the implementation of correct treatment measures. This is true particularly in the case of the above portable unit when used undocked, in which case there is no communication with other equipment.\nFurthermore, U.S. Pat. No. 5,441,047 describes a system according to which selected data on the patient is collected automatically, whereupon said data are forwarded via a stationary telecommunication network, such as a cable television network, to a centre where the diagnosis, monitoring or similar operations may be performed. The referred-to equipment is not, however, portable and in addition it comprises a plurality of independent components, and consequently this equipment is not adapted for ambulatory use and positioning onboard e.g. an ambulance. Nor is it adapted for active exchange of information between the nursing staff by the patient's side and the personnel at the central unit.\nIn addition, there is a need for message exchanges between the nursing staff by the patient's side and the personnel at the central unit as well as for possibilities of filling in certain types of pre-defined forms, such as patient case record files, already in the initial stage by the patient's side. These needs are not met in the prior-art devices.\nIn addition, it would be desirable to enter data manually and preferably by one hand only, in a convenient, rapid and simple manner. It likewise would be desirable, to construct the entry means sufficiently small so as not to make the portable equipment unnecessarily bulky and unmanageable."}
{"text": "Determining the electrical integrity of complex and high density conductive patterns formed on both the exterior and inner layers of ceramic substrates, printed circuit boards and other like items before populating the substrates and boards with components, has become of increased concern to the electronics industry due to rapidly changing technology and economic considerations. Foremost, PC boards in ever shrinking dimensions, are being designed with increasingly dense patterns to accept surface mounting techniques (SMT) of components. These are in the form of extremely reliable, but expensive, chips, DIP's and EPROM's which reduce the need for after assembly \"in-circuit\" and \"functional circuit\" testing, provided the unpopulated boards and substrates are free of defects and leakage paths.\nPrinted circuit board and substrate designs are changed frequently and the lot size of production runs of these items have decreased, but the variety and varying shapes have increased. In board manufacture, the trend is away from the use of expensive metals which have superior bonding characteristics, such as gold and silver, to the use of less reliable metals such as copper, which has a propensity to cause cracks or \"open\" conditions. Lastly, the decision \"not to test\" unpopulated printed circuit boards and substrates can lead to costly scrappage losses after almost error-free component loading under conditions where board continuity failures occur or already existed prior to the mounting of the components.\nIn the past, test equipment to verify circuit integrity and resistivity of PCBs and substrates used either fixtures of multiple probe (bed-of-nails), spring loaded probe cards, or two programmable probes (flying probes). Another known technique is the use of a single probe or capacitance measurement which requires a second connection to the ground plane. All of these known methods have shortcomings. Multiple probe fixtures and cards are fragile, expensive to build, and cause delays since they have to be specially designed for each PCB or substrate design. Dual probe and capacitance testers exhibit relatively slow test cycles in testing dense patterns on substrates and PC boards. Further, probes can make dents or marks on the metallic conductor pads or points to which they are applied, which can cause component bonding failures and may preclude use of surface mounting techniques. Pressure may cause micro cracks to close. Also, high density printed circuit boards of less than 20 mil pad sizes and spacing often preclude the use of \"bed-of-nails\" and probe devices due to high cost factors and physical size constraints. Even a more critical limitation to the use of these known test methods is the possibility of failure to detect micro or hair line cracks in both inner and exposed upper of lower layers. An isolation or \"open\" condition (which occurs where hair line cracks exist) may go undetected by these known systems since they apply an electric potential through the conductors under test. A hair line crack of the order of 1 micron (0.0001 inch) can be bridged by a test voltage as low as 6-10 volts.\nTo overcome many of these problems, an electro-luminescent gas glow discharge continuity testing system was introduced in the early 1980's by Testamatic Corporation for testing the electrical integrity of unpopulated multi-layer ceramic substrates and printed circuit boards. This electro-luminescent gas discharge continuity testing system is described more fully in U.S. Pat. No. 4,507,605 issued March 26, 1985 for a \"Method and Apparatus for Electrical and Optical Inspection and Testing of Unpopulated Printed Circuit Boards and Other Like Items\", assigned to Testamatic, Inc. (the assignee of the present invention), the disclosure of which hereby is incorporated in it's entirety into the disclosure of this application."}
{"text": "There are examples, in which in place of conventional oil lubricating bearings, gas bearings are used as bearings for journal shafts of gas turbines, compressors, etc., which rotate at high speed. In the case where working temperature is high as in gas turbines, thermal expansion and thermal deformation are increased to cause a problem that hydrodynamic journal bearings, of which bearing surfaces opposed to a journal shaft are composed of rigid members and which include a tilting pad and a herringbone groove type hydrodynamic bearings, are hard in terms of reliability to be put to practical use. Therefore, there have been developed and adopted foil hydrodynamic journal bearings, of which bearing surfaces are composed of an elastically compliant foil sheet to be able to absorb thermal deformation, thermal expansion, manufacturing error, and assembly error of respective parts, etc.\nIn addition to the feature described above, foil hydrodynamic journal bearings are stable in high-speed rotation since they elastically support and damp a journal shaft. For example, the literature (TRIBOLOGY TRANSACTIONS, Vol. 43-4 (2000), pages 795–801) describes the historical development of foil hydrodynamic gas bearings.\nAccording to the disclosure of the literature, foil hydrodynamic journal bearings are generally composed of bearing sliding surfaces and spring members, which elastically support the bearing sliding surfaces, and enhanced in bearing performance by an increase in flexibility and damping quality of the bearing. Such foil hydrodynamic journal bearings are mainly classified into leaf type ones and bump type ones according to the structure of spring members, which elastically support bearing sliding surfaces.\nLeaf type foil hydrodynamic journal bearings comprise a plurality of leaf foils arranged round a journal shaft with a bearing gap therebetween in a manner to partially overlap one another on a bearing circle, and a non-cylindrical bearing shape obtained by overlapping of the foils is made use of to generate dynamic pressure upon rotation of the shaft. Also, as an axle load acts, the leaf foils are deformed in a manner to follow the journal shaft, and the bearing is enhanced in damping capacity by friction, which is caused by the deformation to act on the leaf foils, so that it is possible to expect an effect of reduction in rotational vibrations of a journal shaft.\nOn the other hand, bump type foil hydrodynamic journal bearings comprise a substantially circular-shaped top foil arranged round a journal shaft with a bearing gap therebetween, the top foil being elastically supported on a bearing holding member (housing) through a bump foil, which is formed to be corrugated in shape. When a load acts on the journal shaft, the top foil is deformed in a manner to follow the journal shaft, which causes strain on the bump foil and causes slip on contact surfaces between the top foil and the bump foil to give bearing damping by virtue of friction. Besides, multiple winding type foil hydrodynamic journal bearings have been also proposed.\nFurthermore, in order to improve the bearing performance for rotational vibrations or the like, an improvement in stiffness and damping characteristics is made by providing an intermediate layer (mid-foil) to make a plurality of layers, as by means of a method of inserting a further foil between a top foil and an elastic support member (for example, a bump foil), and by contriving the shape and material of the elastic support member. However, foil hydrodynamic journal bearings having the above feature involve a problem that parts are large in number, precision machining is demanded in manufacture, and since handling is difficult to make, a minute work is made necessary at the time of mounting the parts in an assembly stage to be unsuitable for mass production.\nSince in case of a lubricant being gaseous, a lubrication film formed between a bearing sliding surface and a journal shaft upon rotation is as small as several microns or less, minute irregularities are generated on the bearing sliding surface due to being discontinuously supported by a bump foil to make it hard to form an ideal wedge film shape, which constitutes a basis for formation of a lubrication film, with the result that the bearing loading capacity is lowered.\nHereupon, as in a hydrodynamic gas journal bearing described in JP-A-9-14262, there has been proposed an arrangement, in which a vibration damping material (vibration-proof rubber, macromolecular gel, or powder) is provided between a cylindrical-shaped foil and a bearing holding member (housing)."}
{"text": "With regard to methods for administration of drugs, oral administration, rectal administration, percutaneous administration and intravenous administration have known. Among them, oral administration is widely used. In case of oral administration, however, administered drug has disadvantages to receive the first-pass effect after absorption and to show a temporally excessive over blood concentration after administration. Further, in oral administration, many adverse effects such as gastrointestinal trouble, vomiting sensation and anorexia have reported. In addition, in the recent aging society, numbers of patients with dysphagia are increased, and easily ingestible preparations are required from the clinical standpoint.\nAs a result, development of preparations for percutaneous administration has positively promoted for the purpose of solving such disadvantages of oral administration and obtaining easily ingestible preparations for patients with safety and repeatedly, and such the improved preparations are launched in the market.\nHowever, percutaneous absorbability of drug in such preparations for percutaneous absorption is still insufficient. In addition, since many drugs have low percutaneous absorbability, development of the preparations for percutaneous absorption is difficult and the objectives for absorption have not successful. Namely, since the normal skin has barrier function for protecting infiltration of foreign matters into the body, base ingredients used in the conventional preparations for percutaneous absorption can not reveal sufficient percutaneous absorption of the combined active ingredient.\nFor that solution, a contrivance for increasing percutaneous absorption of drug through the stratum corneum of skin is required, and in general so called a promoter for percutaneous absorption is combined with the base ingredient. For example, a combination of absorption promoter of lower alkylamide such as dimethyl acetamide and ethyl alcohol, isopropyl alcohol or isopropyl palmitate (U.S. Pat. No. 3,472,931); a combination of 2-pyrrolidone and proper oil or that of straight chain fatty acid and alcohol ester (U.S. Pat. No. 4,017,641); and a combination of lower alcohol and C7-20 alcohol, C5-30 aliphatic cyclic hydrocarbon, C19-26 aliphatic carboxylic acid alcohol ester, C10-24 mono- or di-ether or C11-15 ketone and water (JP-A-61-249934) have proposed. However, these conventional absorption promoters and the composition of absorption promoter are not sufficiently safe on the skin. Further, in case of the preparations for percutaneous absorption containing salt of basic drug, its effect can not be expected.\nExamples of preparations for percutaneous administration such as a combination of drug and organic acid have reported. For example, tapes combining betamethasone valerate and organic acid in natural rubber adhesives (JP-B-63-45368), tapes combining steroid antiinflammatory agent and organic acid in acrylic adhesives (JP-B-7-47535), and cataplasms combining methyl salicylate as an active ingredient, emulsifier, organic acid, plasticizer, adhesive resin and water in polymer of styrene-isoprene-styrene block copolymer (JP-B-3-31685) have proposed. Objectives of these organic acids are, however, improvement of stability, improvement of solubility and pH adjuster. Further, since these drugs are acidic or neutral, objectives of these preparations are not to improve skin permeability of bioactive substance through formation of ion pair, which is an object to use organic acids in the present invention.\nMeans for improving skin permeability of basic bioactive substances are known. For example, tapes combining citric acid and isoproterenol hydrochloride in acrylic adhesives (JP-A-63-79820), and tapes combining organic acid and vinpocetine in acrylic adhesives (JP-A-5-25039) have reported. These preparations have problems such as irritation in detaching and no sufficient effect for treatment in release of drugs.\nMeans for combining drug and organic acid have reported as preparations for percutaneous administration. For example, a preparation containing organic acid and glycol in nonsteroidal antiinflammatory agent (JP-A-62-181226), and a patch containing alkaline metal salt of nonsteroidal antiinflammatory agent and free nonsteroidal antiinflammatory agent as well as stronger organic acid (JP-B-7-47535) have reported. These inventions relate to acidic drugs and not to basic drugs. Examples of combination of basic drug or its salt, C2-5 alcohol, C2-5 organic acid and C16-20 carboxylic acid ester have known, but no use of organic acid salt is described. In WO 96/16642, patch preparation with a salt of basic drug containing organic acid salt is disclosed, but the fact that organic acid salt increases solubility of basic drug salt in the liquid component having specific solubility parameter ranging from 7 to 13 (cal/cm3)1/2 is not disclosed. In the prior known inventions, preparations for percutaneous absorption, in which basic drug in the form of salt is dissolved in a liquid component to achieve percutaneous absorption of drug, have not known."}
{"text": "1. Field of the Invention\nThe invention relates to a device for creating optical delays and its application to an optical control system for a scanning antenna. More particularly this application relates to a system for generating physical delays for the optical control of scanning antennas. It is a matter of simplifying the structure of the system generating the optical delays and of limiting the number of active elements allowing these delays to be selected.\n2. Discussion of the Background\nSystems are known which allow the control, by optical means, of the radiation pattern of an antenna. They allow control of the phase of the microwave frequency signal transmitted by each antenna element and whose frequency is included in the frequency band 1-12 GHz. The proposed device in particular provides the angular scanning function of the beam on the basis of the control of the phase of each element of the antenna.\nThe present scanning techniques are electronic and act on the respective phases of the controls of the various radiating elements. The phase-shifting elements are ferrite or PIN diode devices. These structures, apart from their weight and consumption, have the disadvantage of only allowing a narrow band functioning of the antenna. In fact, the phase shifters only act between 0 and 2.pi. and therefore, during a change in frequency in the transmitted signal, there is off-aiming of the antenna.\nIt is known, particularly from the French Patent Application N.degree. 87 05267 filed on the 14th Apr. 1987, to produce microwave frequency delay lines with programmable lengths and of small size by optical means. This consists in generating a microwave frequency signal by the beating of two collinear light waves, in making these two light waves travel a path whose length is adjustable by controlled reflectivity electro-optical modulators and in detecting the microwave frequency beat of the two collinear light waves at the end of the path by means of a photodetector.\nAccording to the prior art, it is therefore a matter of encoding the microwave frequency signal, necessary for the control of the active elements of the antenna, on an optical carrier and to physically produce the delays which will allow the phase shifting of this microwave frequency in order to perform the electronic scanning appropriate to this type of radar. In order to produce P phase shifts, it is necessary to generate P optical paths having different lengths and to choose these delays either a priori or a posteriori.\nThe main problem encountered with the solutions proposed previously is the necessity of disposing in the system a number of optical switches (or light modulators) which is equal to the number of delays desired that is, each optical switch. In the case in which the delays are chosen a priori, this results in a high complexity of the system, in a high cost and a possible large attenuation of the optical wave (due to the different modulators). The same problems are also encountered when the choice of the delays is made a posteriori and is even more critical with regard to the attenuation of the optical wave.\nThe purpose of the invention is to provide an encoding system limiting the complexity of the system (but having an equivalent number of available delays as provided by prior systems,) while allowing a reduction in the attenuation of the optical wave and achieving these purposes at a lower cost."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device having a silicon-on-insulator (SOI) structure and a method for fabricating the same, and more particularly, to a method for fabricating a semiconductor device wherein a side wall oxide film or polysilicon layer is formed on the edge of an insulating film adapted to insulate the side wall of an active semiconductor substrate of the semiconductor device from a gate oxide film of the semiconductor device.\n2. Description of the Prior Art\nIn the fabrication of semiconductor devices, formation of a SOI structure is involved to achieve an isolation between adjacent elements, thereby obtaining a superior electrical characteristic. Such a SO structure is made by forming a silicon oxide film as an insulator on an under silicon substrate and forming another silicon substrate (to be used as an active substrate), for example, a single-crystalline silicon layer on the silicon oxide film.\nReferring to FIGS. 1A and 1B, a conventional SOI structure is illustrated.\nFIG. 1A is a view showing the layout of a metal oxide silicon field effect transistor (MOSFET) having the SOI structure. In FIG. 1A, respective positions of masks for an active region 1 and a gate electrode 2 on a silicon substrate (not shown) are shown.\nFIG. 1B is a cross-sectional view taken along the line I--I of FIG. 1A. A silicon oxide film 4 is deposited over a silicon substrate 3. A second silicon substrate 5 having a trapezoidal cross-sectional structure is formed on the silicon oxide film 4. On a desired portion of the resulting structure, a gate oxide film 6 and a gate electrode 7 are sequentially formed.\nAs shown in FIG. 1B, the second silicon substrate 5 has an inclined structure in that its thickness at the edge of the active region is smaller than its thickness d.sub.si at the middle portion of the active region. Due to such a structure, the depth of the depletion region is limited to the thickness of the second silicon substrate 5. As a result, the charge Q.sub.B of depleted bulk is limited by the thickness of the second silicon substrate 5, thereby decreasing the threshold voltage of the semiconductor device, as expressed the following equations: EQU V.sub.T =V.sub.FB +Q.sub.B /C.sub.OX (1) EQU C.sub.OX =.epsilon..sub.OX /t.sub.OX (2)\nwhere, V.sub.T represents the threshold, V.sub.FB the flat band voltage, Q.sub.B the bulk charge, C.sub.OX the capacitance of the oxide film, .epsilon..sub.OX the dielectric constant of the oxide film, and t.sub.OX the thickness of the gate oxide film.\nEquation (1) shows that it is impossible to increase the threshold voltage at the lower edge portion 8 of the second silicon substrate 5 even when the impurity doping rate at that portion increases. This implies a difference in threshold voltage between the middle portion and edge portion of the second silicon substrate 5. The drain current characteristic thus involves a point of inflexion when it varies in accordance with a variation in gate voltage, as shown in FIG. 2. As a result, it is difficult to control the threshold voltage of the semiconductor device having the SOI structure. Furthermore, the characteristic of the SOI element depends on the edge shape of the second silicon substrate 5. This results in a great variation in the characteristics of the semiconductor device."}
{"text": "1. Field of the Invention\nThe invention relates to a method for manufacturing a bending sensor, in particular a strip-shaped bending sensor for an impact sensor device of a pedestrian protection system of a vehicle.\n2. Description of the Related Art\nStudies have shown that a high proportion of fatal accidents in road traffic involve pedestrians. For this reason, legislation initiatives are under way aimed at making the provision of devices for protecting pedestrians in the case of a collision with a vehicle compulsory in modern motor vehicles.\nThere is a particularly high risk of injury for a pedestrian when there is a collision between a vehicle and the distance between a typically easily deformable engine hood and a rigid engine block is very small. The arrangement of increasingly large numbers of electronic components in the region of the engine compartment and the fact that vehicles are of very compact design results in the engine compartment being tightly packed with very rigid elements. If a collision occurs with a pedestrian, there is therefore a high risk of severe head injuries if the pedestrian's head impacts against the engine hood and therefore comes into contact with the components located underneath the engine hood.\nOn the other hand, the risk of injury can be reduced greatly given a sufficiently large distance of, for example, over 10 cm between the engine hood and the engine components arranged underneath it, since the engine hood can absorb a sufficiently large amount of energy as a result of the deformation and can therefore brake the pedestrian's movement comparatively gently.\nIn order to increase safety of pedestrians in road traffic, the European Automobile Manufacturers' Association (ACEA), for example, has given a commitment to the authorities of the European Union to halve the number of traffic deaths involving pedestrians by the year 2010 by means of measures in vehicles. One measure for this purpose is to construct vehicles with correspondingly spaced engine hoods. However, this is frequently impossible due to the required compactness of vehicles.\nIn order to ensure sufficient attenuation in the event of a collision with a pedestrian, it has been proposed, if an impact of a person against the vehicle is detected, to raise the engine hood by more than 10 cm from its closed position in order to provide a sufficient area for deformation. The need for such safety systems to be not only reliable but also very cost-effective presents an enormous challenge for these systems.\nFor example, an article in the specialist journal “Automotive Engineer”, April 2004, page 48 ff., discloses, as an actuator for raising the engine hood, a spring-based actuator whose spring is prestressed and in the event of a detected collision is released, with the result that the engine hood is correspondingly raised. However, pyrotechnic actuators are also known from the abovementioned article.\nCA 2 424 708 A1 discloses a method and a device for detecting a collision between a vehicle and an object. Lightguide fibers are arranged along a front bumper of the vehicle. The lightguide fibers comprise, in their fiber jacket, light exiting regions which are arranged along the lightguide fibers. A collision causes the lightguide fibers to bend. The attenuation of the light which is transmitted in the lightguide fibers changes as a result of the bending of the lightguide fiber if the lightguide fiber is bent in the light exiting region. A signal is acquired from the light which is modulated in this way and said signal is processed in a signal processor. A safety device, for example for raising an engine hood, can be activated in this way."}
{"text": "FIG. 1 shows a portion of a semiconductor memory comprising memory cells MC arranged in a so called “virtual ground array.” The gates G of memory cells MC arranged along rows are connected by the same wordline WL. The source/drain regions SD of memory cells MC arranged along columns are connected to the same local bitlines LB, with each local bitline LB being shared by memory cells MC in two adjacent columns of the array in order to reduce chip area. Each local bitline LB is connected to a total of 512 memory cells MC while each wordline WL is connected to the gates G of 512 memory cells MC. The memory cells MC are nitride read-only memory (NROM) cells, which are non-volatile. NROM cells can store two bits per cell in a nitride layer. The ability to store two bits per cell is indicated in the schematic of each memory cell MC by two “x”. The combination of NROM cells with virtual ground arrays allows the design of memories with high storage densities.\nIn order to select memory cells MC for reading, each local bitline LB can be connected to one of the global bitlines GB by means of a respective switching element S1. The global bitlines GB are connected to column decoders (not shown) and the wordlines WL to row decoders (not shown). Potentials applied to the global bitline GB selected by the column decoder and to the wordline WL selected by the row decoder are then passed on to the source/drain regions SD and the gate G of the memory cell MC to be read. The state stored in the memory cell MC can be determined by sensing the current flowing through the memory cell MC in a sense amplifier.\nTo minimize the chip area required for the semiconductor memory, the local bitlines LB are alternately connected to top and bottom global bitlines GB, with six local bitlines LB being connectable to each global bitline GB. For ready reference, a switching unit SU comprising switching elements S1, part of a global bitline GB and parts of local bitlines LB is shown.\nFIG. 2 shows the layout of a switching unit SU. Six local bitlines LB are connected by means of local bitline contacts CL to diffusion zones in a semiconductor substrate SB. Further shown is a global bitline GB which is connected to further diffusion zones in the semiconductor substrate SB by means of global bitline contacts CG. Global bitline select lines SG are arranged between each local bitline contact CL and the global bitline contact CG, which is closest to the respective local bitline contact CL.\nEach combination of a local bitline contact CL, its closest global bitline contact CG and the global bitline select line SG between the two contacts forms a transistor. The gate of the transistor is connected by the global bitline select line SG and the source/drain regions of the transistor are connected by the local bitline contact CL and the global bitline contact CG, respectively. If a potential greater than the threshold potential of the transistor formed is applied to the global bitline select line SG, then the local bitline LB connected by the local bitline contact CL is connected to the global bitline GB connected by the global bitline contact CG. The transistor is thus an implementation of one of the switching elements S1 shown in FIG. 1 and, because it is used to select a bitline, is commonly known as select transistor.\nThe technology that is used to manufacture the semiconductor memory determines the smallest pitch, that is, the smallest distance between lines, that is possible. The size of the memory cells MC determines the pitch of the layout, and as each memory cell MC is directly connected by the respective local bitlines LB, this pitch is also the pitch for the local bitlines LB.\nThe local bitlines LB cannot be directly connected to the sense amplifier as they are too small and too resistive. Therefore, another metallization layer with global bitlines GB is introduced. The global bitlines GB are thick metal layers of low resistance and are used to connect the local bitlines to the sense amplifier. However, because of the larger pitch of the global bitlines GB, it is not possible to connect all of the local bitlines LB to different global bitlines GB. Rather, each of the global bitlines GB has to be shared by six local bitlines LB, so that it is not possible to control each of the local bitlines LB individually without using complicated architectures or decoding mechanisms.\nFIG. 3 is used to illustrate an effect known as “neighbor effect” which occurs when a local bitline LB is shared between two memory cells connected by the same wordline WL. For the sake of clarity, only one row having three memory cells is shown. However, it is clear that the cells can be extended to the left and right as well as to the top and bottom as shown in FIG. 1. In addition to the elements already described in conjunction with FIG. 1, a sense amplifier SA connected to the first global bitline GB1 and to the second global bitline GB2, is shown.\nIn FIG. 3, memory cell MC is selected for reading by closing the switching elements S1 of the local bitlines LB connecting the source S and drain D of the memory cell MC. As a result, these local bitlines LB are connected to the first global bitline GB1 and the second global bitline GB2, respectively. The local bitlines LB of the neighboring memory cell NC and the further memory cell FC, which are not shared with the memory cell MC, are not connected to the global bitlines GB1 and GB2. A current IS will flow through the memory cell MC if suitable potentials VS and VD are applied to the first global bitline GB1 and to the second global bitline GB2, respectively. An erased memory cell MC allows a higher current IS to flow than a programmed cell, so that the state stored in the memory cell MC can be determined by measuring the current IS flowing through it.\nIdeally, the current IM flowing into the sense amplifier SA is equal to the current IS flowing through the memory cell MC. However, if a neighboring cell NC is connected to the same local bitline LB as the memory cell MC, some of the current IS flowing through the memory cell MC will leak through the neighboring memory cell NC. This current IN will depend on the state stored in the neighboring cell MC and may flow even when no potential is applied to the gate of the neighboring cell NC. Current may, therefore, leak through all the memory cells that are connected to the same local bitline LB, so that the total leakage current maybe as large as 10 to 30% of the current IS flowing through the memory cell MC.\nIn FIG. 3, the current IS flowing out of the source S of the memory cell MC is to be measured. This is referred to as source-side sensing. It is also possible to measure the current ID flowing out of the drain D of the memory cell MC. This is known as drain-side sensing and the current flowing into the second global bitline GB2 is measured. If a current IF flows through the further memory cell FC leakage will also occur in drain-side sensing, adulterating the current measured.\nAs a consequence of the leakage due to the neighbor effect, the current IM measured in the sense amplifier SA is less than the current IS or ID flowing through the memory cell MC. If the leakage current is great enough, then the current IM measured may be decreased to such an extent that a programmed memory cell MC is mistakenly read as an erased cell. This will lead to reading failure of the memory as data cannot be correctly retrieved. There is, therefore, a need to reduce the leakage due to the neighbor effect as much as possible.\nIn prior art, the neighbor effect problem has been partially solved by charging or discharging the local bitlines LB and the global bitline GB before each read operation.\nIn source-side sensing, all bitlines, global and local, are pre-discharged to the same value as the potential VS applied to the first global bitline GB1 for reading, which typically is 0 V. As a result, the voltage across the source and drain of neighboring cells is approximately zero so that no current flows through the neighboring cells. For drain-side sensing, all the bitlines, global and local, are pre-charged to the same value as the potential VD applied to the second global bitline GB2 for reading, which typically is the supply potential. The switching elements S1 connecting the local bitlines LB of memory cells, which are not to be read, are opened after the pre-charging or pre-discharging.\nHowever, when reading the memory cell MC, voltage drops of 100 mV to 300 mV occur, depending on the state stored in the memory cell MC, so that the neighboring cells also display a potential difference of the same magnitude across source and drain. As a result, a small leakage current will still flow. With reductions in structure sizes the resistance in the path for the current IS will increase and larger voltage drops will occur, leading to more leakage. The pre-charge/pre-discharge solution of the prior art, therefore, fails to completely solve the problems of the neighbor effect.\nBesides failing to totally solve the leakage problem, the charge/discharge operations increase the power consumption of the semiconductor memory as all the global and local bitlines must be charged or discharged. A further disadvantage is that the time required for reading the memory cells is increased as the charge/discharge operation must be performed before each read operation and due to the RC time constants involved this takes a certain time.\nPrior art also suggests connecting further global bitlines GB at different potentials to the neighboring cells in order to reduce the neighbor effect. However, providing biased potentials is difficult and these solutions usually require complicated interleaved bitline architectures with large increases in the chip area in order to be able to control the potentials of the local bitlines LB independently from each other. Additionally, the global bitlines GB must still be charged and discharged before each read operation, thus increasing the power consumption of the memory."}
{"text": "Olive oil extraction is the process of extracting the oil present in the olive drupes for food use. The oil is produced in the mesocarp cells, and stored in a particular type of vacuole called a lipovacuole. Every cell contains a tiny olive oil droplet. Olive oil extraction is the process of separating the oil from the other fruit contents. After washing the olives different known techniques may be employed for extracting oil. Traditional methods typically press the olive using a grinder with two millstones. First the olives are ground into an olive paste using large millstones. The olive paste is generally ground by the stones for 30 to 40 minutes for guaranteeing that the olives are well ground and for allowing the fruit enzymes to produce some of the oil aromas and taste. After grinding, the olive paste is spread on fiber disks, which are stacked on top of each other, then placed into a press. These disks are then put on a hydraulic piston, forming a pile. Pressure is applied on the disks, thus compacting the solid phase of the olive paste and percolating the liquid phases. To facilitate separation of the liquid phases, water flows on the sides of the disks to increase the speed of percolation. The liquids are then separated either by a standard process of decantation by gravity or by means of a faster vertical centrifuge. The traditional method provides good results and high quality of oil, although the machine is rather difficult to clean. It is a non continuous process with waiting periods thus exposing the olive paste to the action of oxygen and light. This method typically requires additional manual labor and a longer time period from harvest to pressing as compared to other methods for producing olive oil.\nModern methods of olive oil extraction use an industrial decanter to separate all the phases by centrifugation. In this method the olives are first crushed to a fine paste. This can be done by a hammer crusher, disc crusher or knife crusher. The paste is then malaxed for 30 to 40 minutes in order to allow the small olive droplets to agglomerate. The aromas are created in these two steps through the action of fruit enzymes. Water is added to facilitate the extraction process with the paste during malaxation. Typically the added water is heated in order to produce higher extraction percentage of olive oil. However, water temperatures above 35 degrees Celsius adversely affect the quality of the produced oil. Afterwards the paste is pumped into an industrial decanter where the phases are separated. The decanter is a large capacity horizontal centrifuge rotating at approximately 3000 rpm. The high centrifugal force created allows the phases to be readily separated according to their different densities (solids>vegetable water>oil). Although this technique is continuous and automatic providing high percentages of oil extraction it still produces large amounts of vegetable water to be disposed of and a reduced amount of antioxidants due to the added water.\nYet another known method for extracting oil, also know as “Sinolea” is a process which employs rows of metal discs or plates which are dipped into the paste. The oil preferentially wets and sticks to the metal and is removed with scrapers in a continuous process. The method is based on the different surface tension of the vegetable water and the oil; these different physical behaviors allow the olive oil to adhere to the steel surface while the other two phases stay behind. However large surface areas can lead to rapid oxidation of the olive product. In addition, the cleaning process of the surfaces is difficult and time consuming.\nAll above methods are designed for large scale industrial press applications and for handling of large quantities of olives and oil.\nThere is hence a long need for a system and method for extracting oil which may also be used as a house hold counter-top appliance. Such desired method may not require any additives nor trained personnel."}
{"text": "Many operating systems (OS) today do not include a built in mechanism, called user exits, to divert program control from the operating system or related services to user supplied functions. In many programming instances it is desirable to divert or intercept system calls, issued from a user application, and execute user supplied code instead. The code supplied by the user might bypass the entire original system function call or it might perform a function in conjunction with the original system function call. One application of such a system might include a security system whereby operating system calls issued by a user application are not permitted to execute unless the calling process has the requisite authority or privileges. Another application of this system of providing user exits might include a system profiler whereby any or all operating system calls issued by an application are tracked and used to generate statistics about the behavior of the user application. Yet another application of the system is a system used for enhancing operating system services where it is not desired to permanently modify the internal kernel functions built into the operating system.\nUser exit are used extensively on mainframe computer. They provide a mechanism to allow a system routine to be executed when a running program arrives at certain points. The system routine, however, might contain code supplied by the user for performing a particular function. For example, user supplied system code might function to determine whether a process should be granted or denied permission to use certain system services or resources. The security example presented above is similar in that user exits are employed to deny or allow user applications access to system services such as communications channels, access to devices or other system resources.\nImplicit in the concept of providing user exits on an operating system platform, is that the supplied user exit code itself is able to be inserted and installed within the operating system where it can be accessed. Some operating systems provide the capability of inserting user supplied code into kernel space. Kernel space refers to the area of memory that contains the core operating system software including associated system calls and various data structures needed to manage and operate the computer platform on which the operating system is installed. Other operating systems, however, do not provide a capability of inserting user code into kernel space."}
{"text": "This invention relates to the manufacture of semiconductor devices as components of integrated circuits, specifically to processes for photolithography in which anti-reflective coatings are used to increase the accuracy of the photolithographic processing steps.\nIn a dual damascene architecture, wiring patterns are etched into a dielectric, or insulating layer. See, e.g., Handbook of semiconductor interconnection technology, edited by Schwartz et al., Marcel Dekker 1998; and Copperxe2x80x94Fundamental Mechanisms for Microelectronic Applications, Murarka et al., Wiley 2000. The conductor material (typically copper) is then inlaid into those features. There are two types of features used for this purpose: trenches, which form the actual wiring template; and vias, which make connection to the metal level below. Creating such structures requires two passes through the photo-lithography process. Either the vias are formed first and then the trenches or vice versa.\nThe dielectric stack requirements for dual damascene include the primary insulating layer and a thin copper diffusion barrier or selective metal barrier. Additional layers may be included to facilitate fabrication, such as an intermediate etch stop, hard mask, etc.\nAn anti-reflective layer (or ARL) is often used for photolithographic processes. The ARL minimizes the total reflection of light from layers under the photoresist and the interface between the photoresist layer and the underlying layer. By adjusting the thickness, t, refractive index, n, and extinction coefficient, k, of the ARL film, as shown in FIG. 1, a destructive interference can be obtained in the photoresist with equivalent intensities of incident and reflective light. As a result, zero reflectivity can be reached under ideal conditions. Thus, an ARL improves the accuracy of pattern transfer when the photoresist is developed.\nFIG. 2 show a schematic diagram of a simplified lithography process flow of the via-first dual damascene applications. Typically, anti-reflective layer 18 is deposited onto underlying layer 12, which is being patterned over the other film stack. Photoresist 10 is then spun onto top of the anti-reflective layer. See, FIG. 2(a). The process proceeds through (b) via photoresist development to (c) via etch to (d) via photoresist removal and cleaning to (e) or (e1) trench photoresist coating and (f) or (f1) trench photoresist development. The exposed portion of the photoresist layer 10 is removed when photoresist layer 10 is developed, yielding the clean vertical walls shown in FIG. 2(b) when UV radiation is incident on area of the top surface of photoresist layer 10, exposing a portion of the photoresist layer 10. When developed, the trench should be patterned properly and yielded vertical wall, as shown in FIG. 2(f).\nHowever, this identity in pattern after the development step is not always realized. More specifically, silicon dioxide (SiO2) historically has been used as the primary interconnect insulating layer. With device geometries shrinking and speeds increasing, the trend now is towards insulating materials with lower dielectric constants (low-k). One of the most persistent difficulties associated with the integration of the low-k film has been its interaction with photoresists used with deep UV radiation (xe2x80x9cDUVxe2x80x9d, i.e., radiation having a wavelength of 248 nm and below). Low-k films often contain a small amount of nitrogen, present in the form of NHx (amines). The NHx species can diffuse rapidly through low-k dielectric films, such as would be used in layer 12. Such groups are known to react in a detrimental fashion with DUV photoresists by neutralizing the photo-acid catalyst. The result is footing or bridging of the printed features. These footings narrow the opening in the photoresist which results in poor pattern transfer to the underlying layers. See, FIG. 3.\nTo alleviate this phenomenon, a hard mask could be incorporated on top of the low-k film. This approach is very costly and is limited to hard masks that are barriers to the diffusion of amines (i.e., low k hard masks would not be suitable.) Moreover, the use of a hard mask will be effective only for single layer lithography; e.g., printing vias for via-first dual damascene, or trenches in the case of trench-first. However, after formation of those features into the dielectric, DUV resist would once again come in contact with low-k film during the next pass through lithography. The result would be regions of undeveloped resist in the second pass lithography features. For example, in the case of via-first patterning, the appearance of xe2x80x9cmushroomsxe2x80x9d or xe2x80x9crivet headsxe2x80x9d over the vias (also filled with resist) in the trench regions. A DUV photoresist xe2x80x9cmushroomxe2x80x9d is illustrated schematically in FIG. 2(f1) and demonstrated in a SEM picture shown in FIG. 4.\nOrganic layers have been commonly used as ARLs for I-line radiation (i.e., radiation having a wavelength of 365 nm), although inorganic layers can also be used. Layers of inorganic materials such as silicon oxynitride are often used for deep UV radiation. For a discussion of the use of deep UV ARLs, see T. Ogawa et al.; xe2x80x9cPractical Resolution Enhancement Effect By New Complete Anti-Reflective Layer In KrF Excimer Laser Lithographyxe2x80x9d; Optical/Laser Microlithography; Session VI; Vol. 1927 (1993), incorporated herein by reference. For a general discussion of the use of ARLs, see T. Perara, xe2x80x9cAnti-Reflective Coatings: An Overviewxe2x80x9d; Solid State Technology, Vol. 37, No. 7; pp. 131-136 (1995), which is incorporated herein by reference.\nThere is a need for a simpler approach to deposit an ARL film with elimination of the photoresist footing problem and photoresist poisoning (mushroom) problem in dual damascene processes. It would be desirable for the ARL film to be optically and thermally stable, to be chemically inert to the environments to which it is exposed, and to be applicable for use with any wavelength of UV radiation. In addition, the ARL film should have good adhesion to commonly used materials and have good mechanical and structural integrity. Finally, it would be desirable to provide a single continuous process for producing the ARL film with acceptable uniformity across the wafer .\nIn accordance with this invention, a substantially nitrogen-free ARL is provided. The ARL is compatible with chemically amplified photoresists found in lithography processes using DUV radiation. It is effective in eliminating footing and inhibits the photoresist poisoning of the second layer of lithography in dual damascene applications.\nThe ARL can be formed using chemical vapor deposition (CVD), including, for example, plasma enhanced chemical vapor deposition (PECVD), low-pressure chemical vapor deposition, (LP-CVD), high-density plasma chemical vapor deposition (HDP-CVD), or similar methods under process conditions that are designed to provide an ARL having the desired properties. These properties principally include thickness (t), refractive index (n), and extinction coefficient (k), with the values of n and k being a function of the wavelength of the radiation.\nThe ARL is formed by introducing source gases or liquids comprising silicon, oxygen, carbon, and hydrogen into the reaction chamber of the CVD unit. In a particularly preferred embodiment, the ARL is formed from carbon dioxide and silane. The process conditions or parameters of this invention are principally the flow rate of the reactant gases, the volumetric ratio of the gaseous components, the deposition pressure, the deposition temperature, and the rate at which radio frequency (RF) power is applied in a PECVD unit (per unit area of the surface of the wafer or other substrate on which the ARL is formed)."}
{"text": "1. Field of the Disclosure\nThe invention relates generally to jack stands for vehicles and specifically to universal jack stands that can be used with a variety of vehicle jacks and jack stands that support a vehicle at a common jacking location with the jack.\n2. Related Technology\nElevating a portion of a vehicle typically involves using some sort of jack. For example, many automobiles carry a scissor-type jack to elevate a portion of the automobile in order to change a tire. The jack is placed under a portion of the frame of the automobile and the jack is slowly raised until a platform on the jack engages the frame of the automobile. Thereafter, a user uses the principle of leverage to elevate a portion of the frame. Often, more than one location on the frame needs to be elevated so that a user can repair a portion of the automobile that is only accessible from the bottom. Because most standard jacks are movable, so that the jack may be easily moved to many different portions of the frame, a vehicle supported solely by the jack may be unstable and unsafe to work under. For this reason, a user may place a jack stand under the frame and lower the jack so that the vehicle is supported on the jack stand instead of the jack itself. The jack stand is typically a stable platform for supporting the vehicle. Because the jack and jack stand have individual and separate support structures, the jack stand and jack typically cannot support the vehicle at same location on the frame. This can cause a problem when a manufacturer designates only certain locations as jacking locations. Generally, vehicle manufacturers reinforce certain locations on the frame of a vehicle so that those certain locations can support a portion of the weight of the vehicle when the vehicle is elevated with a jack (hereinafter these locations on the vehicle frame are referred to as “jacking locations”).\nGenerally, jack stands are pyramid-shaped structures having three or four sides and an adjustable support platform disposed in the top of the jack stand. For example, U.S. Pat. No. 7,207,548 discloses a more or less typical jack stand including four sides that form the shape of a pyramid, a telescoping support member extending from a top portion of the pyramid and a curved horizontal support, often having a U-shape, at the top of the telescoping support member to cradle a portion of the vehicle frame. The support member is often adjustable with a ratchet type mechanism. Problems occur with placement of this type of platform as only certain parts of the frame or undercarriage can fit in the curved horizontal support. While a jack stand is usually more stable than a jack, the jack stand cannot support the vehicle at the same point that is occupied by the jack. As a result, the jack stand often must be placed a considerable distance away from the jacking location in order to find a portion of the frame or undercarriage that is compatible with the top of the support platform, necessarily requiring a higher jacking elevation to accommodate the jack stand location.\nIn order to solve this problem, systems have been designed to elevate the vehicle with the jack stand itself. For example, U.S. Patent Publication No. 2008/00999745 discloses a three stage jack stand. A lift collar on the top of the jack stand is elevated by a specially designed unique power unit. The power unit includes a pair of lift arms that engage a lower portion of the lift collar on the jack stand. When the jack stand is placed under a vehicle, the lift arms on the power unit move upward to engage a bottom of the lift collar on the jack stand. The lift collar fully supports and elevates the vehicle at all times during the jacking process. One drawback to this type of arrangement is that the jack stand and the power unit are specially designed for one another. Thus, this type of jack stand and power unit are not compatible with contemporary vehicle jacks, such as scissor jacks, floor jacks or racing jacks. Moreover, this type of power unit is not compatible with typical jack stands and/or may be too tall to fit under some vehicles."}
{"text": "Portable communication devices, such as cellular telephones, personal digital assistants (PDAs), WIFI transceivers, and other communication devices transmit and receive communication signal at various frequencies that correspond to different communication bands and at varying power levels. A power amplifier module, generally comprising one or more amplification stages, is used to transmit the communication signals. A radio frequency (RF) power amplifier system may include multiple amplification stages, and, in some applications, multiple amplification paths. The efficiency of the power amplifier system is generally determined by a number of factors, and to a large degree, determines the amount of power consumed by the power amplifier system.\nVarious ways of measuring and determining the efficiency of the power amplifier are available. One measure of power amplifier efficiency is referred to as “power added efficiency,” abbreviated as PAE. The PAE of a power amplifier is dependent upon a number of factors including, but not limited to, the number of power amplification paths, the load impedance at the output of the power amplifier path or paths, impedance matching between multiple stages, and other factors. In a multiple mode power amplification topology, in which two or more power amplification paths are implemented to provide varying power output levels, improving the PAE at low power levels is typically achieved at the expense of power amplifier linearity at higher power levels.\nTherefore, it is desirable to improve the PAE of a power amplifier system over a range of power levels, without sacrificing the linearity and performance of the power amplifier system over the range of power output levels."}
{"text": "Hydrolyzable resins typified by polylactic acid are excellent in biodegradability and are now under study as substitutes for various plastics in various application fields from the viewpoint of environmental improvement, and some of them have been put to practical use.\nIn recent years, use of the resins as additives to be added to drilling fluids used to extract underground resources has been proposed (refer to Patent Document 1).\nFor example, a winze drilling method called “hydraulic fracturing method” is now widely employed to extract underground resources. In this drilling method, a high pressure is applied to a drilling fluid filled in a winze to form cracks (fractures) in the vicinity of the winze so as to improve permeability (ease of the fluid's flow) in the vicinity of the winze and to expand the effective sectional area through which a resource such as oil or gas flows into the winze, thereby increasing the productivity of the winze. This drilling fluid is also called “fracturing fluid”, and a viscous fluid such as gel-like gasoline was used. However, an aqueous dispersion prepared by dissolving or dispersing polymer particles in water has recently been used due to the development of shale gas produced from shale strata existent at relatively shallow sites in consideration of an influence on the environment. A hydrolyzable resin such as polylactic acid is proposed as this polymer.\nThat is, polylactic acid is a substance which exhibits hydrolyzabibility and biodegradability and is decomposed by water or an enzyme contained in the ground even if it remains in the ground and therefore does not exert a bad influence on the environment. Water used as a dispersion medium can be considered to have almost no influence on the environment as compared with gasoline.\nWhen a drilling fluid containing hydrolyzable resin particles such as polylactic acid dispersed in water is filled in a winze and pressurized, the resin particles permeate an area in the vicinity of the winze and become a sealing material for the formed fractures to temporarily block a flow channel of a resource such as gas or oil effectively. To form new fractures, the fluid is also used as a diverting agent for switching a flow channel of a fracturing fluid. Further, since the fluid is hydrolyzed in the winze and disappears, the hydrolyzable particles do not need to be removed in a post-step, thereby making it possible to carry out the drilling of the winze efficiently.\nBy the way, a large number of studies have been made on the heat resistance and mechanical properties of a hydrolyzable resin such as polylactic acid but almost no studies are made on its particle shape.\nFor example, Patent Document 2 discloses a biodegradable resin composition (polylactic acid composition) comprising polylactic acid having a D-isomer content of 2 mass % or less and lamellar silicate. However, Patent Document 2 is aimed to improve the heat resistance and mechanical properties of this polylactic acid composition and does not investigate its particle shape at all.\nPatent Document 3 discloses a powder comprising polylactic acid having a crystallinity of 30% or more. This polylactic acid powder has such high mechanical grindability that it can be particulated. However, Patent Document 3 does not investigate its particle shape as well.\nFurther, the inventors of the present invention proposed a polylactic acid composition which comprises low crystalline or amorphous polylactic acid having a crystallinity of 40% or less as a matrix and has a dispersion structure that an organically modified polysaccharide or swollen or expanded layered silicate salt is dispersed as a grinding accelerator in the matrix (Japanese Patent Application No. 2014-47835).\nThis polylactic acid composition has advantages that it has excellent mechanical grindability and is easily particulated at a low cost to be usable as an additive for drilling fluids. That is, since an agent to be added to a liquid, such as a drilling fluid which is used in a large quantity is used in a large quantity as a matter of course, the possibility of preparing it as fine particles which are easy to be injected into a liquid such as water and easy to handle at the time of injection at a low cost brings about a great benefit in terms of cost.\nAlthough the mechanical grindability of particles used as an additive for drilling fluids is studied in the above previous application, the particle shape is not investigated at all as well."}
{"text": "1. Field of the Invention\nThis invention relates to an optical type scanning microscope. This invention particularly relates to a scanning microscope, wherein an optical means, which irradiates a light beam to a sample, is moved with respect to a sample supporting member, on which the sample is supported, such that the light beam may scan the sample. This invention also relates to a device for detecting a scanning width, i.e., a width over which a light beam, or the like, scans a material, in an apparatus in which the material is scanned with the light beam, or the like. This invention additionally relates to a magnification indicating apparatus for a scanning microscope, such as an optical type scanning microscope or a scanning electron microscope, wherein an actual width of scanning with a probe is detected, and a magnification of the microscope is indicated in accordance with the detected scanning width.\n2. Description of the Prior Art\nOptical type scanning microscopes have heretofore been used. With the scanning microscope, a light beam is converged to a small light spot on a sample, and the sample is two-dimensionally scanned with the light spot. Light radiated out of the sample during the scanning (i.e., light which has passed through the sample, light which has been reflected from the sample, the fluorescence which is produced by the sample, or the like) is detected by a photodetector. An enlarged image of the sample is thereby obtained. An example of the scanning microscope is disclosed in Japanese Unexamined Patent Publication No. 62(1987)-217218.\nIn the conventional optical type scanning microscopes, a mechanism which two-dimensionally deflects a light beam by a light deflector is primarily employed as the scanning mechanism.\nHowever, the scanning mechanism described above has the drawback in that a light deflector, such as a galvanometer mirror or an acousto-optic light deflector (hereinafter referred to as \"AOD\"), which is expensive must be used. Also, with the scanning mechanism described above, a light beam is deflected by a light deflector. As a result, the angle of incidence of the deflected light beam upon an objective lens of the light projecting optical means changes momentarily, and aberration is caused to occur. Therefore, the scanning mechanism described above also has the problem in that it is difficult for the objective lens to be designed such that aberration can be eliminated. Particularly, in cases where an AOD is utilized, astigmatism occurs in the light beam radiated out of the AOD. Therefore, in such cases, a special correction lens must be used, and the optical means cannot be kept simple.\nIn order to eliminate the aforesaid problems, a scanning mechanism has heretofore been proposed wherein a light beam is not deflected but a sample is scanned with the light spot of the light beam. For example, in U.S. Pat. No. 5,081,350, a novel mechanism has been proposed wherein a light projecting optical means is supported on a movable member, the movable member is moved reciprocally with respect to a sample supporting member, and a light spot of a light beam is thereby caused to scan a sample.\nIn U.S. Patent application Ser. No. 735,734, the applicant has proposed a movement mechanism for moving an optical means with respect to a sample supporting member. The proposed movement mechanism comprises a tuning fork having an end, on which an optical means or a sample supporting member is supported, and an excitation means composed of an electromagnet for applying a magnetic field, the intensity of which changes periodically, to the tuning fork, and thereby causing the tuning fork to vibrate. The proposed movement mechanism is advantageous over a movement mechanism utilizing, for example, a piezo-electric device or an ultrasonic vibrator, in that the width, over which the optical means moves with respect to the sample supporting member, i.e. the width, over which the light beam scans the sample, can be kept large. (The width, over which the light beam scans the sample, is determined by the amplitude of the tuning fork.) Therefore, when the proposed movement mechanism is employed in a scanning microscope, an image of a large area of the sample can be formed.\nIn most of conventional scanning microscopes, light radiated out of a sample (i.e., light which has passed through the sample, light which has been reflected from the sample, the fluorescence which is produced by the sample, or the like) is detected by a photodetector, and a serial output generated by the photodetector is fed into a signal processing means. The signal processing means samples and quantizes the output received from the photodetector. In this manner, image signal components of a digital image signal are obtained which correspond to respective positions lying along each main scanning line on the sample.\nHowever, in such cases, if a light beam scanning mechanism, which is composed of a tuning fork and an electromagnet, is utilized, problems may occur in that a distortion occurs in an image which is reproduced from the digital image signal. This is because the displacement of the tuning fork does not change linearly with respect to time but changes in accordance with characteristics close to a sine function. Such problems also occur when a light beam scanning mechanism other than the light beam scanning mechanism composed of a tuning fork and an electromagnet is utilized.\nIn order to cope with the problems described above, it is considered to modulate the period of pixel clock pulses that determines the period, with which the output of a photodetector is sampled, so as to compensate for fluctuations in the speed, at which an optical means is moved reciprocally with respect to a sample supporting member. Such a method is effective, but it has the problem in that, if the relationship between the operation of the tuning fork (i.e., the movement of the optical means with respect to the sample supporting member) and the timing, with which the sampling process is begun, is not accurately kept at predetermined relationship, a distortion occurs in an image reproduced from the digital image signal.\nThe aforesaid scanning mechanism, wherein a light beam is not deflected but a sample is scanned with the light spot of the light beam, is broadly applicable to optical type scanning microscopes, scanning electron microscopes, scanning tunnel microscopes, light beam scanning read-out apparatuses, or the like. In such a scanning mechanism, a movement mechanism composed of a tuning fork and an electromagnet for resonating the tuning fork, a movement mechanism composed of a piezoelectric device, a movement mechanism composed of an ultrasonic vibrator, or the like, is often utilized in order to move a scanning means (e.g., a probe, on which a light projecting optical means, a stylus, or the like is supported, in cases where the scanning mechanism is employed in a scanning microscope) reciprocally and quickly with respect to a material which is to be scanned. In cases where such a movement mechanism is employed, the width, over which the scanning means moves reciprocally with respect to the material to be scanned, changes comparatively easily due to fluctuations in the drive voltage applied to the movement mechanism, or the like. If the width, over which the scanning means moves reciprocally with respect to the material to be scanned, changes, the width over which the light beam scans the material will change inevitably.\nIn cases where the scanning mechanism is employed in a scanning microscope, if the width over which the light beam scans the material changes, the magnification, with which a microscope image is formed, will deviate from a designed value. In cases where the scanning mechanism is employed in a light beam scanning read-out apparatus, if the width over which the light beam scans the material changes, the size of an image, which has been read out from the scanned material, will fluctuate. In order to eliminate such problems, it is necessary that the actual width, over which the light beam scans the material, is detected by a certain means. In accordance with the detected scanning width, the magnification of the microscope image should be indicated, or the operation of the scanning mechanism should be corrected such that the scanning width may be kept at a predetermined value.\nAs one of scanning width detecting devices for detecting the actual width, over which a light beam scans a material, there has heretofore been known a device comprising:\ni) a grid pattern, which is combined with either one of a material to be scanned and a scanning means and which is constituted of a plurality of light reflecting members or light blocking members standing side by side with one another in a direction, along which the scanning means is moved reciprocally with respect to the material to be scanned, PA0 ii) a light projector, which is combined with the other of the material to be scanned and the scanning means and which irradiates a light beam to the grid pattern, PA0 iii) a light receiver for detecting the light beam, which has been reflected by the grid pattern, or the light beam, which has passed through the grid pattern, and PA0 iv) a means for counting the number of periodical fluctuations in a light beam detection signal, which has been generated by the light receiver. PA0 i) a sample supporting member on which a sample is supported, PA0 ii) an optical means which irradiates a light beam to said sample, PA0 iii) a movement mechanism which reciprocally moves said optical means with respect to said sample supporting member such that said light beam may scan said sample in main scanning directions and in sub-scanning directions, and PA0 iv) a photodetector for detecting light radiated out of the portion of said sample, which is exposed to said light beam, an image of said sample being thereby formed, PA0 a) a signal processing means for sampling a serial output of said photodetector and thereby generating image signal components of a digital image signal, which correspond to each main scanning line, PA0 b) a means for feeding pixel clock pulses into said signal processing means, the period of said pixel clock pulses being modulated so as to compensate for fluctuations in a speed, at which said optical means is moved reciprocally with respect to said sample supporting member, PA0 c) a means for generating a timing signal that determines the timing, with which the sampling process is begun, PA0 d) a grid pattern, which is secured to either one of said sample supporting member and said optical means and which is constituted of a plurality of light reflecting members or light blocking members standing side by side with one another in the direction, along which said optical means is moved reciprocally with respect to said sample supporting member, PA0 e) a light projector, which is associated with the other of said sample supporting member and said optical means and which irradiates a light beam to said grid pattern, PA0 f) a light receiver for detecting the light beam, which has been reflected by said grid pattern, or the light beam, which has passed through said grid pattern, PA0 g) a means for sampling a light beam detection signal, which has been generated by said light receiver, in accordance with said pixel clock pulses and thereby generating a digital displacement signal, and PA0 h) a means for approximating the relationship between the order x, where x=1, 2, 3, . . . , in which signal components of said digital displacement signal representing specific points appearing at equal pitches on said grid pattern occur, and the order y, in which said signal components of said digital displacement signal are sampled, with a quadratic equation y=ax.sup.2 +bx+c, PA0 i) a material, which is to be scanned, PA0 ii) a scanning means, and PA0 iii) a movement mechanism which reciprocally moves said scanning means with respect to said material such that said scanning means may linearly scan said material, PA0 a) a grid pattern, which is combined with either one of said material and said scanning means and which is constituted of a plurality of light reflecting members or light blocking members standing side by side with one another at predetermined pitches in a direction, along which said scanning means is moved reciprocally with respect to said material, PA0 b) a light projector, which is combined with the other of said material and said scanning means and which irradiates a light beam to said grid pattern, PA0 c) a light receiver for detecting the light beam, which has been reflected by said grid pattern, or the light beam, which has passed through said grid pattern, PA0 d) a means for sampling a light beam detection signal, which has been generated by said light receiver, in accordance with predetermined sampling clock pulses and thereby generating a digital displacement signal, and PA0 e) a calculation means for approximating the relationship between the order x, where x=1, 2, 3, . . . , in which signal components of said digital displacement signal representing specific points appearing at equal pitches on said grid pattern occur, and the order y, in which said signal components of said digital displacement signal are sampled, with a simple equation y=ax+b, PA0 i) a sample, PA0 ii) a probe, and PA0 iii) a movement mechanism which reciprocally moves said probe with respect to said sample such that said probe may linearly scan said sample, PA0 a) the aforesaid scanning width detecting device in accordance with the present invention, PA0 b) a means for calculating a magnification of an image reproduction width in a microscope image reproducing means, which image reproduction width is taken in the direction of said scanning, with respect to the scanning width, which has been found by said scanning width detecting device, and PA0 c) an indicating means for indicating the magnification, which has thus been calculated.\nIn the scanning width detecting device\nfluctuations in described above, the number of periodical the light beam detection signal represents the number of the light reflecting members or the light blocking members, across which the light beam irradiated from the light projector to the grid pattern has moved. Therefore, the width, over which the light beam scans the material, can be detected by counting the number of periodical fluctuations in the light beam detection signal. For example, in cases where the grid pattern is constituted of a plurality of the light reflecting members, and the light receiver detects the light beam, which has been reflected by the light reflecting members, the light beam detection signal generated by the light receiver periodically rises pulse-wise. Therefore, the actual width, over which the light beam scans the material, can be detected by counting the number, N, of the pulse-wise rising points in the light beam detection signal and multiplying the pitch, p, at which the light reflecting members of the grid pattern stand side by side with one another, by the number, N.\nHowever, with the scanning width detecting device described above, errors easily occur in detecting the width, over which the light beam scans the material. Therefore, the scanning width detecting device described above cannot always be employed when the width, over which the light beam scans the material, is to be detected very accurately. How such errors occur will be described hereinbelow. By way of example, the number, N, is counted as being 100 in cases where the light beam irradiated from the light projector to the grid pattern has moved across 100 light reflecting members or 100 light blocking members (i.e., m'th through m+99'th members). Actually, in some cases, the light beam may begin scanning immediately before it is irradiated to the m'th member, and the scanning may be finished immediately after the light beam has moved across the m+99'th member. In other cases, the light beam may begin scanning when it is irradiated to a position in the vicinity of an m-1'th member, and the scanning may be finished at a position in the vicinity of an m+100'th member. In both cases, the width, over which the light beam scans the material, is detected as being equal to 100p. However, in the former cases, the actual width, over which the light beam scans the material, is close to 99p. In the latter cases, the actual width, over which the light beam scans the material, is close to 101p.\nIn order to eliminate the problems described above, it is considered to set the initial position, at which the scanning begins, at a predetermined position. However, in cases where a piezo-electric device, or the like, is employed as the reciprocal movement mechanism, it is very difficult to set the initial position, at which the scanning begins, at a predetermined position."}
{"text": "1) Field of the Invention\nThe present invention relates to a transmission characteristics evaluation system, and more particularly to a transmission characteristics evaluation system which is suitable, for example, when used in evaluating the characteristics of an optical communication transmission system.\n2) Description of the Related Art\nIt is known in the art that, together with the communication performance of light that is emitted by a laser diode (LD) used in an optical communication system, the signal conversion performances of an electric/optical (E/O) conversion module for converting an electric signal into an optical signal and an optical/electric (O/E) conversion module for converting an optical signal into an electric signal are dependent on the dispersion characteristics of an optical fiber constituting a transmission path.\nParticularly, the dispersion tolerance, which is a transmission characteristic of the optical fiber constituting the aforementioned transmission path relative to the magnitude of the amount of dispersion, is a critical index value in evaluating the performance quality of an optical transmitting apparatus provided with an LD or an E/O module or an optical receiving apparatus provided with an O/E module described above.\nFIG. 13 is a block diagram illustrating a conventional measurement system for measuring the dispersion tolerance. In this measurement system 200 shown in FIG. 13, an optical transmitting apparatus 210 and an optical receiving apparatus 220 are connected with a dummy fiber 230 which is assumed to be an optical fiber constituting the transmission path in a real optical communication system.\nHere, the optical transmitting apparatus 210 is provided with a pattern pulse generator (PPG) 204 together with an E/O module 203 comprising a light source 201 and a modulator 202, whereby an electric signal of a specific pulse pattern generated by the pattern pulse generator 204 of the E/O module 203 is converted into an optical pulse signal by the modulator 202 and is transmitted to the optical receiving apparatus 220 via the aforementioned dummy fiber 230.\nFurther, the optical receiving apparatus 220 is provided with a receiving amplifier 221, an O/E module 222, and a BER tester (Bit Error Rate tester) 223, whereby the optical pulse signal transmitted via the dummy fiber 230 (and amplified by the receiving amplifier 221) is converted into an electric signal by the O/E module 222, and is output to the BER tester 223 as a received electric signal.\nThe BER tester 223 performs error detection on the received electric signal relative to the electric signal of the pulse pattern generated by the pulse pattern generator 204 of the optical transmitting apparatus 210. When the dispersion value (fiber length) of the dummy fiber 230 used here is changed, the error rate detected by the BER tester 223 changes. By using this change in the error rate relative to the change in the fiber length, the dispersion tolerance of the optical element, the optical module, and the measurement system can be measured.\nIn other words, with this error detection value obtained from the BER tester 223, the communication quality of the E/O module 210 and the O/E module 220 relative to the dispersion value of the dummy fiber 230 can be evaluated. In addition, by increasing or decreasing the dispersion value of this dummy fiber 230, the communication quality of the E/O module 210 and the O/E module 220 can be measured as a dispersion tolerance.\nHere, in measuring the error detection value by increasing or decreasing the dispersion value of the dummy fiber 230 in the above-described measurement system 200, it is necessary to connect a dummy fiber 230 having a different length newly to the optical transmitting apparatus 210 and the optical receiving apparatus 220 each time the error detection value is measured with the dummy fiber 230.\nRecently, there is known a measurement technique that uses a variable dispersion compensator of a Fiber Bragg Grating (FBG) type (optical network simulator manufactured by JDS-Uniphese (US), TeraXion (Canada), or SPIRENT Co., Ltd.) instead of this dummy fiber. These techniques eliminate the need for replacing the dummy fiber 230 in order to increase or decrease the dispersion value of the transmission path, such as in the above-described measurement system 200.\nOn the other hand, in recent years, in evaluating the transmission characteristics of an optical transmitting apparatus provided with an LD or an E/O module or an optical receiving apparatus provided with an O/E module described above, the insertion loss gradient tolerance, which is the amount of change in the insertion loss of the optical fiber constituting the aforementioned transmission path in accordance with the wavelength of the transmitted light, is considered as one of the index values.\nAs a technique for measuring this insertion loss gradient tolerance, a variable wavelength bandpass filter module having a bandwidth about ten times larger than the original bandwidth for use is interposed between the optical transmitting apparatus 210 and the optical receiving apparatus 220 such as described above and, by shifting the center wavelength of this filter module, a gradient is given to the distribution of the insertion loss relative to the wavelength of the used transmittance band, and the signal error rate is measured with the BER tester 223.\nNamely, the insertion loss gradient tolerance is measured by measuring the signal error rate that changes in accordance with the degree of gradient of the distribution of the insertion loss relative to the wavelength of the used transmittance band.\nHere, as a technique related to the present invention, a technique disclosed in the following patent document 1 is known.                [Patent Document 1] Japanese Patent Application Laid-open No. 2002-258207        \nHowever, in the case of measuring an index value for evaluating the performance quality of an optical module such as an optical transmitting apparatus provided with an LD and an E/O module or an optical receiving apparatus provided with an O/E module, a technique is demanded having a high measurement precision with reduced number of working steps for the measurement while widening the measurement range as compared with such a conventional technique.\nFor example, the aforementioned conventional technique for measuring the dispersion tolerance raises a problem in that, in order to perform evaluation of multiple channels, i.e. in order to evaluate the dispersion tolerance of each wavelength in transmitting a wavelength multiplex light, a variable dispersion compensator interposed between the optical transmitting apparatus and the optical receiving apparatus must be prepared separately in accordance with the wavelength of the object of measurement. Moreover, the dispersion value cannot be changed arbitrarily from a positive dispersion value to a negative dispersion value, so that the dispersion tolerance for evaluating the transmission performance of the optical module cannot be measured in a sufficient range.\nAlso, the aforementioned conventional technique for the measurement of insertion loss gradient tolerance raises a problem in that the insertion loss gradient tolerance cannot be measured commonly in the aforementioned measurement system for measuring the dispersion tolerance. Moreover, when the center wavelength of the bandpass filter module is shifted, a change occurs not only in the value of insertion loss but also in the dispersion value within the band. This makes it difficult to evaluate correctly whether the error detection value detected by the BER tester 223 is due to the change in the dispersion value or due to the change in the value of insertion loss."}
{"text": "This disclosure relates generally to a deaerator and, more particularly, to a deaerator having an apertured stem that facilitates separating air from a coolant.\nGenerators provide electric power as is known. Aircraft auxiliary power units, for example, typically include a generator that is driven by a turbine engine. The turbine engine is rotatably coupled to the generator through a gearbox. Coolant, such as oil, is circulated through the gearbox and the generator. The coolant removes thermal energy and lubricates various components. Coolant mixes with air as the coolant circulates through the generator. As known, at least some of the air must be separated from the coolant before the coolant can be reintroduced to the gearbox and the generator. Deaerators are used to separate air from the coolant. The coolant is collected within a sump after the deaerator removes the air. The coolant is recirculated through the gearbox and the generator from the sump."}
{"text": "The present invention relates to an image sensing apparatus suitable for using a sequential scanning type solid-state image sensing device having a color filter of so-called Bayer arrangement, and the like.\nRecently, an image sensing device, such as CCD, capable of sequentially reading signals of all the pixels (referred as xe2x80x9cnon-interlace scanning type image sensing devicexe2x80x9d, hereinafter) has been developed with the progress of semiconductor manufacturing technique. The non-interlace scanning type image sensing device has an advantage in that a higher resolution image can be obtained with less blurring than an image sensed by using a conventional interlace scanning type image sensing device even when sensing a moving object. In the interlace scanning type image sensing device, a frame image is composed of two field images which are sensed at different times, usually at a field period interval. Accordingly, there is a problem in which, when sensing a moving object, there are notches on edges of the object and perhaps of the background in a frame image because of the time gap between the two field images composing a frame image. If a frame image is made of image data of a single field image to overcome the aforesaid problem, there would not be notches on edges, however, the vertical resolution of the obtained frame image is halved. In contrast, with a non-interlace scanning type image sensing device, it is possible to sense a frame image in a field period, thus, the non-interlace scanning type image sensing device is expected to be applied to a still image camera and a camera for EDTVII, for example.\nA general image sensing apparatus using an image sensing device which outputs signals after adding two vertically adjacent pixel charges will be explained with reference to FIG. 22.\nReferring to FIG. 22, an image sensing device 901 output signals after adding two vertically adjacent pixel charges in accordance with timing signals t1 and t2 generated by a timing signal generator (TG) 909. The output image signals are inputted to a correlated double sampling (CDS) circuit 903 via a buffer 902, and reset noises of the image sensing device 901 are removed from the output image signals by the CDS circuit 903., then the image signals enter an automatic gain controller (AGC) 904. In the AGC 904, the image signals are amplified by a gain set in accordance with a control signal c2 from a microcomputer 908 (gain control). The gain-controlled image signals are converted into digital signals by an analog-digital (A/D) converter 905, then transmitted to a camera processing circuit 906, where predetermined processes are applied to the digital image signals and a luminance signal Y and a color difference signal C are outputted. Further, the microcomputer 908 generates a control signal c2 for controlling the gain in the AGC 904 in accordance with the gain information c1 detected by a camera processing circuit 906.\nMost of the non-interlace scanning type image sensing devices used at the present time are provided with R, G and B filter chips arranged in so-called Bayer arrangement as shown in FIG. 23. In this color arrangement, G signal is used as a luminance signal. Further, the non-interlace scanning type image sensing devices output image data of one frame in a field period, thus the speed for transferring charges is two times faster than the transferring speed of an image sensing device, as shown in FIG. 22, which outputs image signals obtained by adding two vertically adjacent pixel charges. Accordingly, it is preferred to design an image sensing device to have two horizontal registers which respectively transfer image signals of odd and even lines simultaneously to be first and second output signals, instead of transferring by one line through a single horizontal register.\nIn this case, the buffer 902, the CDS circuit 903, the AGC 904, and the A/D converter 905 shown in FIG. 22 are needed for each of the two horizontal registers. In a case where image signals are obtained from an image sensing device provided with a color filter of Bayer arrangement as shown in FIG. 23, the G signals are obtained from all the corresponding pixels by horizontally interpolating by using signals obtained at G filter chips (i.e., G signals) as shown in FIG. 24A. The G signals obtained as above are converted into luminance signals to be displayed. However, if image signals to be displayed are obtained in this manner, difference in characteristic between the two AGCs will affect the quality of an image.\nThe AGCs change gains to be applied to image signals when a gain is provided, and the two AGCs have different characteristics from each other in general. Therefore, even though the same gain is provided to the two AGCs, the levels of amplified image signals may differ from each other. If this occurs, when an object of a uniform color (e.g., a white paper) is sensed, a variation in output signal level of the two AGCs appears in a stripe pattern of odd and even lines. Therefore, when an image of the object is displayed, the variation in output signal level appears as the difference in output signal level between odd and even line fields on a display as shown in FIG. 24B, which causes field flicker. This noticeably deteriorates the quality of the image.\nTo overcome this problem, a method for interpolating an average of pixel values of the G signals in a vertical row can be considered. However, in this method, when an object of a uniform color (e.g., a white paper) is sensed, a variation in output signal level of the two AGCs may cause vertical stripes which alternatively have different output levels on a display as shown in FIG. 24C. The difference in output level in the vertical stripes also noticeably deteriorates the quality of an image.\nFurther, in order to compensate a variation in output signal level of the two AGCs, feed-back control is considered. Feed-back control can be performed in the following sequence, for example. First, a test signal is inputted to the two horizontal registers of the image sensing device at predetermined timing, then the test signals outputted from the two horizontal registers are processed in the same manner of processing image signals. Then, the difference in output level between the two AGCs is detected. On the basis of the detected difference in output level, new gains to be provided to the AGCs are adjusted. Thereafter, a test signal is inputted to the horizontal registers again so as to confirm that the difference in output level between two AGCs are corrected.\nAs described above, if there is a variation in characteristic between two AGCs, it is possible to compensate a variation in output signal level between the two AGCs by performing feed-back control.\nHowever, in the aforesaid method, a problem in which the size of a hardware increases is posed in providing a feed-back control function for compensating a variation in output signal level between two AGCs, since a circuit necessary for performing the feed-back control has to be added.\nFurther, because of an arrangement of pixels for luminance signals (i.e., pixels with xe2x80x9cGxe2x80x9d filter chips) as shown in FIG. 23, the following problem arises. FIG. 25 is a graph showing a spatial frequency plane. In FIG. 25, the abscissa denotes a spatial frequency in the horizontal direction, the ordinate denotes a spatial frequency in the vertical direction, and the oblique axis denotes the spatial frequency in the oblique direction. In FIG. 25, the spatial frequencies at ∘ marks are the spatial sampling frequency decided in accordance with the number of pixels, and so on, of CCD (the spatial sampling frequencies in the horizontal, vertical and oblique directions are respectively referred as xe2x80x9cfs-hxe2x80x9d, xe2x80x9cfs-vxe2x80x9d, and xe2x80x9cfs-oxe2x80x9d, hereinafter). When signal components having these sampling frequencies are sampled, the obtained signals appear as a direct current component signal, which gives bad effects on an image. Therefore, it is necessary to remove signal components of the sampling frequencies from the image.\nAs a method for removing such the frequency components from the image, an optical crystal low-pass filter (referred as xe2x80x9cO-LPFxe2x80x9d, hereinafter) is generally used. More specifically, with an O-LPF designed to have null points at spatial frequencies shown by three ∘ in FIG. 25 in the oblique direction, it is possible to effectively remove sampling frequency components not only in the oblique direction but also in the horizontal and vertical directions.\nHowever, moire due to signal components, shown by ⊚ in FIG. 25, having Nyquist frequencies which are half of the sampling frequencies (i.e., xc2xd(fs-h), xc2xd(fs-v), and xc2xd(fs-o)) also appear on an image, which deteriorates the image. Especially, the moirxc3xa9 which appears in the oblique direction has very large component signals, and it appears even more clearly when displaying a moving image and an image obtained while panning a camera. However, the aforesaid O-LPF can not remove such moirxc3xa9.\nIn order to overcome this problem, it is possible to use an O-LPF designed to have null points at spatial frequencies shown by ⊚ in FIG. 25 in the oblique direction, similarly to the O-LPF for removing signal components of sampling frequencies, to remove signal components of Nyquist frequencies in the horizontal, vertical and oblique directions. However, in this methods, a bandwidth of luminance signals obtained when an image is formed on the non-interlace type image sensing device becomes narrow. Furthermore, when another filter is used, response in the bandwidth drops dramatically. Thus, a resultant image will not be expressed sharply.\nIn addition, an image sensed by the aforesaid non-interlace scanning type image sensing device can be displayed both by a non-interlace method and by an interlace method. Accordingly, as an image sensing apparatus having a non-interlace scanning type image sensing device, one which can switch between outputting an image for a interlace type display and outputting an image for a non-interlace type display is demanded.\nFurther, there is a demand for inputting photograph data to a personal computer. Especially, there is a demand for taking an image into the personal computer from a negative film since images sensed by using a camera are often stored in negative films.\nThe present invention has been made in consideration of above situation, and has as its object to effectively prevent deterioration of the quality of an image due to a variation in amplification characteristic in each line upon reading image signals by a plurality of lines and due to an arrangement of a color filter.\nFurther, it is another object to prevent deterioration of the quality of an image due to a variation in amplification characteristic in each line upon reading image signals by a plurality of lines and due to an arrangement of a color filter by using a simple circuit configuration without adding a feed-back type correction circuit.\nAccording to a preferred configuration of the present invention, each color signals obtained by separating image signals which are read from an image sensing device are processed by a properly designed two-dimensional filter. With this two-dimensional filter, it is possible to remove signal components of Nyquist frequencies in the horizontal, vertical and oblique directions, thereby harmful signals due to a variation in amplification characteristic (gain), which appears as a variation of luminance signals having Nyquist frequencies, corresponding to each of a plurality of horizontal registers of the image sensing device and due to an arrangement of a color filter are effectively removed. As a result, the quality of an image is improved.\nIt is still another object of the present invention to enable to switch between image signal outputs conforming to a plurality of display methods in an image sensing apparatus.\nFurther, it is still another object of the present invention to enable to change characteristics of image signals to be outputted depending upon the state of an object and the type of output device. More specifically, it is an object to provide an image sensing apparatus capable of setting characteristics of edge enhancement signal frequencies in accordance with a state of an output device.\nFurthermore, it is still another object of the present invention to reduce the size of a circuit for negative-positive inversion as well as to provide an image signal control circuit for achieving a sufficient white balance adjustment.\nFurther, it is still another object of the present invention to realize a circuit which reads and processes image signals from a non-interlace scanning type image sensing device by inputting two alternate scan lines so as to be able to switch between different signal output formats (e.g., interlace output and non-interlace output) with a simple circuit configuration."}
{"text": "At present, there exist a variety of cutting machines for cutting metal members, among which the gantry-type numerical control cutting machine is widely used. This machine comprises a cross beam, a driving trolley, and a driven trolley, wherein the driving trolley and the driven trolley are symmetrically and rigidly connected to the bottom surface of the welded cross beam at two ends thereof, at two ends thereof, said rigid connection being in such a way that a connecting plate is connected to the bottom surface of the cross beam, and the connecting plate is connected to the driving trolley and the driven trolley. Such a machine has the following shortcomings and deficiencies:\n1. The cross beam is individually welded. The manufacturing is difficult, costly, and it is not easy to ensure manufacturing precision. Meanwhile, the manufactured beam is big, heavy, and has an unpleasant appearance.\n2. As the two ends of the cross beam are rigidly connected to the driving trolley and the driven trolley, stress and deformation may be incurred due to certain reasons during the operation of the cutting machine. This affects the stable operation of the cutting machine, and in serious circumstances machine parts may be damaged."}
{"text": "1. Field of the Invention\nThe aspects of the present disclosure relate generally to LED light sources, and in particular to an LED light source compatible with different phase control dimmers.\n2. Description of Related Art\nMost dimmers used in lighting applications, such as home lighting, commercial lighting, and traffic signaling, are designed for use on locally available AC grid voltage. These dimmers operate with higher power consumption lamps that generally present as resistive loads to their controllers. Recent advancements in lamp technology have led to development of Light Emitting Diode (LED) lamps with improved efficiency and significantly lower power consumption. These LED lamps typically require DC power and often consume less than ten watts. Thus these lamps are not directly compatible with existing lighting fixtures and dimmers.\nNearly all lamp dimmers designed for use with a local AC power grid include a triac device to reduce lamp power. A triac is a bidirectional thyristor device for alternating current which may be viewed as a switch that can conduct current in either direction. When these dimmers are used to drive LED lamp assemblies, flickering can occur due to the switching and high current peaks resulting from triac switching current fluctuations caused by interaction with an EMI input filter. This flickering can have detrimental effects on the triac device. Further, a typical LED lamp is not compatible with all types of phase controlled dimmers. The lack of universal compatibility will prevent or drastically reduce the introduction of LED lamps into the marketplace as replacements or retrofits for standard incandescent lamps.\nIncreasing the power consumption of LED lamps is one approach that has been used to adapt LED lamps to various LED lamp dimmers. However, increasing the power consumption can create thermal issues in the LED lamp and the associated circuitry and offsets many of the benefits that drive the move to low power LED lamps. For example, the increased temperatures can reduce the usable life of the LED lamp, and increased power consumption reduces cost benefits.\nOther approaches for adapting LED lamps to various dimmers have included the use of digital integrated circuit solutions that detect a firing angle of the dimmers or triacs. However, LED lamps do not always present as a resistive load leading to inaccurate detection of firing angle. Although the use of specific drivers, specific integrated circuits or more expensive digital integrated circuits may provide certain solutions, these solutions are typically directed to specific LED lamps and will work with only a few, specific dimmers. It would be advantageous to be able to utilize LED lamps in a wide range of existing and new applications and with existing and new dimmers without concern for compatibility issues. It would also be advantageous to maintain the power factor of a LED lamp circuit near unity, when a LED lamp is directly connected to the local mains supply without a dimmer.\nAccordingly, it would be desirable to provide LED lamps that solve at least some of the problems identified above."}
{"text": "Polypropylene films have found wide acceptance in the packaging industry, especially as a food packaging agent, because of their superior physical properties. Polypropylene film, usually biaxially oriented, is characterized by high tensile modulus and stiffness, and excellent optical clarity and a certain degree of moisture resistance. However, polypropylene film is highly pervious to gases and air. Moreover, polypropylene possesses one major disadvantageous property, high inherent coefficient of friction (COF).\nHigh COF complicates the processing of polypropylene film. Polypropylene film processing is impeded by poor transport, caused by its high COF, over rollers, guides and the like. In addition, high COF created film storage problems. Because of its high COF one layer of polypropylene film sticks to those above and beneath it, creating the problem commonly referred to as blocking.\nThis serious disadvantage of polypropylene film is well known to those skilled in the art. Thus, many proposed solutions to overcome this major deficiency have been proposed in the art. One such proposal has been to incorporate additives in the polypropylene resin processed to form the film. A favorite additive suggested in the art is fatty acid amides. These amides decrease the film's COF as they migrate to the film surface after heat treatment and aging.\nAlthough this method has been used, actual COF is a function of the heat history to which the film has been exposed during shipping, storage and processing. As such, it is subject to wide variation. More significantly, the presence of fatty acid amides on the film surface oftentimes adversely affects the appearance of the film as manifested by decreased gloss and the presence of streaks. Another serious disadvantage of using fatty acid amides is the detrimental effect of fatty acid amides on polypropylene film surface wettability and adhesion. This adverse characteristic applies to coating, inks, adhesives and the like, especially in water based forms.\nIt is also known to coat polypropylene film with certain fatty acid amides to impart lubricating and antiblocking characteristics, as taught in U.S. Pat. No. 4,255,644. However, the application of such coatings by the film manufacturer is not particularly attractive because of the requirement that they be applied as solutions in organic solvents. Health and safety factors dictate against the in-plant utilization of organic solvents in coating processes undertaken during the film manufacture.\nOther slip additives have been suggested to overcome the inherent problem of high COF in thermoplastic films. One such solution, as disclosed by U.S. Pat. No. 4,302,506, is the use of a latex coating containing stearamidopropyl- dimethyl-beta-hydroxyethylammonium nitrate and a crosslinkable acrylic copolymer. This coating has been advanced for use on polyester films.\nAlthough this coating is alleged to improve slip properties without adversely affecting clarity in polyester films it is unnecessary to determine whether this coating would provide the same improved result on polypropylene film. It is known that such a coating is not in conformity with the requirements and regulations of the U.S. Food and Drug Administration for use in films which contact foods. As such a critical market for polypropylene films, as a food packaging agent, could not be exploited if such a material was added to polypropylene-containing film.\nYet another suggestion advanced in the prior art to reduce the blocking characteristic of polypropylene film is to add a finely divided inorganic material as a surface modifier. Such a modification not only provides a non-blocking surface having improved slip characteristics but is also independent of the heat history of the film. Moreover, the addition of such a surface modifier does not create the adverse optical and wettability effects associated with amide-modified films.\nHowever, like the other proposed solutions, this proposed means of alleviating high COF creates new problems. Oftentimes, polypropylene film is laminated to other films. For example, glassine paper is commonly laminated to polypropylene film. Such a laminate, when provided as a thin surface layer containing finely divided inorganic material, exhibits significantly higher COF values than does the unlaminated polypropylene film. Such laminates are known to perform marginally at best on conventional form, fill and seal machines.\nThe above discussion reflects the need in the art for a new and improved polypropylene film characterized by improved anti-blocking characteristics and decreased COF. However, it is incumbent that the improvement in film slip property not correspond to the decline in other properties, typical of the solutions advanced in the prior art."}
{"text": "Today, the use of interactive collaboration systems is becoming increasingly pervasive. Today's systems make it much easier for users to collaborate. However, the uses of automated attendants that are part of a communication session have been limited. These automated attendants mainly act in the background and are not dynamic or act outside of a communication session. For example, a leader of a conference may record the conference by entering a number such as *7. However, there is no mechanism where these types of attendants can be added dynamically using a graphical user interface. Moreover, these attendants are not able to actively take the place of a user of the communication session. For example, if the user cannot make a conference call, the user may have to try and contact the other participants or send a representative to the conference. What is needed is a way to overcome the limited capabilities of the current automated attendants and provide users with increased flexibility for managing communications."}
{"text": "Clark and Gollan's discovery triggered the investigation on the possible use of PFC as oxygen-carrying fluids and Slovitor and Geyer's breakthrough using PFC in an emulsified form has made it realistic. One of the most important factors relating to the realization of such oxygen-carrying fluids is the synthesis of PFC which are non-toxic, have high oxygen and carbon dioxide solubilities, are stable as an emulsion for 2 years even at room temperature; and are easily removable from the body in an unchanged form after accomplishing its role with the formation of the natural blood.\nVarious PFC are disclosed in patents, wherein they are described as being suitable as an oxygen and carbon dioxide carrier. Almost all of the PFC patented are prepared by the electrochemical fluorination method or the cobalt trifluoride method. The cobalt trifluoride method, which is excellent in preparing perfluorohydrocarbon, is too vigorous to fluorinate ethers. That is, degradation products are caused by the cleavage of oxygen and carbon bonds. The electrochemical fluorination method, which is excellent in preparing perfluoro amine, has limited utility in preparing perfluoro cyclic ethers. Therefore, in general, perfluorocyclic ethers, which are believed to be suitable as an oxygen and carbon dioxide carrier, remain to be synthesized. The exception are the ones synthesized by the electrochemical fluorination method, such as those set forth below: ##STR1## described by Abe et al (e.g., T. Abe, E. Hayashi, H. Baba and S. Nagase: J. Fluorine Chem. 25:419 (1984), and Japanese Patent Application (OPI) Nos. 119449/79, 119471/79, 128566/79 and 44071/80.\nIn the present invention, the substrates followed by the prefix \"perfluoro\" mean substrates having all of the hydrogens replaced by fluorine atoms but containing no halogens such as chlorine, bromine and iodine, and they may be cyclic or straight-chain compounds. Furthermore, polyfluorinated compounds with a few chlorine atoms denote substrates wherein the hydrogens are mostly replaced by fluorines, but the rest of the hydrogens, usually one or two hydrogens, are replaced by chlorine atoms."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of consumer packaging, and more specifically to metal cans, such as the steel and aluminum cans that are commonly used for packaging soft drinks, other beverages, food and aerosol products.\n2. Description of the Prior Art and Recent Technology\nMetal cans for soft drinks, other beverages and other materials are of course in wide use in North America and throughout the world. The assignee of this invention, Crown Cork & Seal Company of Philadelphia, is the world's largest designer and manufacturer of such cans.\nThe art of making and packing metal cans is constantly evolving in response to improved technology, new materials, and improved manufacturing techniques. Other forces driving the evolution of technology in this area include raw material prices, the nature of new materials to be packaged and the marketing goals of the large companies that manufacture and distribute consumer products such as soft drinks.\nInterest has existed for some time for a metal container that is shaped differently than the standard cylindrical can in such a distinctive way to become part of the product's trade dress, or to be otherwise indicative of the source or the nature of the product. To the inventors best knowledge, however, no one has yet developed a practical technique for manufacturing such an irregularly shaped can at the volume and speed that would be required to actually introduce such a product into the marketplace.\nU.S. Pat. No. 3,224,239 to Hansson, which dates from the mid 1960's, discloses a system and process for using pneumatic pressure to reshape cans. This process utilized a piston to force compressed air into a can that is positioned within a mold. The compressed air caused the can wall to flow plastically until it assumed the shape of the mold.\nTechnology such as that disclosed in the Hansson patent has never, to the knowledge of the inventors, been employed with any success for the reshaping of drawn and wall ironed cans. One reason for this is that the stress that is developed in the wall of the can as it is being deformed can lead to defects that are potentially failureinducing, e.g., localized thinning, splitting or cracking. The risk of thinning can be reduced by increasing the wall thickness of the can, but this would make shaped cans so produced prohibitively expensive. The risk of splitting and cracking can be reduced by a process such as annealing, but at the expense of reduced toughness and abuse resistance of the final product.\nA need exists for an improved apparatus and process for manufacturing a shaped metal can design, that is effective, efficient and inexpensive, especially when compared to technology that has been heretofore developed for such purposes, and that reduces the tendency of a shaped can to fail as a result of thinning, splitting or cracking."}
{"text": "The present invention relates to a device for delivering ink, more particularly to a print head for an ink jet printer, provided with a series of ink delivery nozzles which can be actuated selectively when printing a sheet of paper or the like. The present invention also relates to an ink jet printer provided with such a device.\nOne known device of this kind has a substantially triangular base plate, the apex of which is directed towards the sheet for printing, and the oblique sides of which are provided with a number of ink ducts covered by a foil and extending to near the bottom or head end of the base plate, where they extend in the form of ink tunnels through the base plate towards the head surface. There the nozzles of the ink tunnels are covered by a nozzle plate, possibly in the form of a foil, in which a hole is made for each ink passage. Piezo-electric elements are disposed at the location of the foil over the ducts on the oblique sides and can be actuated selectively in order to briefly press on the foil above a duct and thus create a pressure pulse in the ink there, so that a droplet is delivered from the associated nozzle.\nU.S. Pat. No. 4,364,067 discloses a print head attachment having a substantially rectangular base plate, the top of which is provided with an ink inlet connected to an ink reservoir and terminating in an ink distribution chamber to which are connected a number of ducts formed in the two sides. These ducts converge with an arc in the plane of the associated side to a number of downwardly directed ink delivery nozzles provided in a nozzle plate disposed in the base plate head surface situated transversely of the sides. The ink delivery nozzles form a Y-shaped duct, the branches of which extend to ducts on the opposite sides of the base plate.\nThe ducts themselves are covered by a vibrating plate on which an electrode is mounted. Provided on the electrode, at the location of the pressure chambers, are piezo-electric elements, which are also provided with an electrode. Selective actuation of the piezo-electric elements causes the vibrating plate to be pressed in locally and the volume of the pressure chamber in the required duct is reduced so that a specific quantity of ink is propelled through said duct to the associated ink delivery nozzle.\nA disadvantage of the known printer heads is that the ink-tight layers adjoin one another at an angle so that sealing may be locally inadequate."}
{"text": "Patients with metastatic renal cell carcinoma have a very poor prognosis with a five-year survival rate of &lt;5%. Recently, biological therapies (e.g., interferons, interleukins, monoclonal antibodies) have demonstrated moderate anti-tumor activity in this disease. In some studies, over 30% of patients treated with alpha-interferons have exhibited a major response. Interferon treatment is associated with moderate toxicity, however, at present it is impossible to predict which subset of patients will benefit from interferon therapy. Monoclonal antibodies, while of limited therapeutic benefit, have been useful in diagnostic studies. Some of these antibodies identify as assortment of glycoproteins of varying molecular weights expressed on the cell surface of renal cancers."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus and method for providing passive shimming for a superconducting magnet of the type that has a relatively small imaging volume. Such structures of this type, generally, provide the proper amount of passive shimming for a relatively small imaging volume.\n2. Description of the Related Art\nIn magnetic resonance (MR) imaging systems, the magnetic field inhomogeneities have to be less than a few parts per million (ppm) for producing good images. Therefore, a superconducting magnet is designed for a homogeneous field distribution of only a few (ppm) within the imaging volume. However, once built, due to the manufacturing tolerances, lead routings, and other manufacturing difficulties, the ppm inhomogeneities are usually in the range of hundreds. Active or passive shimming needs to be used to homogenize the magnetic field.\nIn recent years, due to its lower cost, simplicity, and capabilities, passive shimming has been the preferred method of field inhomogeneity correction. Exemplary of such passive shimming systems are U.S. Pat. Nos. 4,698,611 and 4,771,244 entitled \"Passive Shimming Assembly for MR Magnet\" and \"Method for Passively Shimming Magnetic Resonance Magnet\", respectively, both to M. E. Vermilyea and both assigned to the same assignee as the present invention. While these prior passive shimming assemblies have met with a modicum of success, due to the compact nature of an imaging system which images only the limbs, the space available for passive shimming is extremely limited. In order to keep the length of the MR system down and the inside diameter of the MR system large, the magnet coils and the gradient coil need to be physically very close. This is in fact the area where passive shims need to be placed. Therefore, a passive shimming systems needs to be designed which will occupy the least amount of physical space.\nIt is apparent from the above that there exists a need in the art for a passive shim/gradient coil assembly which is light weight and compact through simplicity of parts and uniqueness of structure, and which at least equals the shimming and gradient coil characteristics of known shimming and gradient coil systems, but which at the same time is capable of providing passive shimming and a gradient coil for a relatively small imaging volume. It is a purpose of this invention to fulfill this and other needs in the art in a manner more apparent to the skilled artisan once given the following disclosure."}
{"text": "In recent years, accompanied with popularization of an optical disk device which has compatibility between a DVD and a CD, a monolithic semiconductor laser is used as a light source, which emits laser light having two wavelengths from places apart from each other with a determined interval by forming an InGaAlP based semiconductor laser element for red light and an AlGaAs based semiconductor laser element for infrared light on the same semiconductor substrate (cf. for example PATENT DOCUMENT 1).\nAs shown in FIG. 3, the AlGaAs based semiconductor laser element 50a is provided, on a semiconductor substrate 51 made of, for example, n-type GaAs, with, for example, an infrared light emitting layer forming portion 59a, which is composed of an n-type cladding layer 52a made of n-type AlGaAs based semiconductor, an active layer 53a made of AlGaAs based semiconductor and a p-type cladding layer 54a made of p-type AlGaAs based semiconductor and having a ridge shape, and a current constriction layer 55a for an infrared element made of n-type GaAs, which is formed on sides of the ridge portion, and a contact layer 56a made of, for example, p-type GaAs is provided thereon. On the other hand, the InGaAlP based semiconductor laser element 50b is provided, on the same semiconductor substrate 51 where the AlGaAs based semiconductor laser element 50a is formed, with, for example, an red light emitting layer forming portion 59b, which is composed of an n-type cladding layer 52b made of n-type InGaAlP based semiconductor, an active layer 53b made of InGaAlP based semiconductor and a p-type cladding layer 54b made of p-type InGaAlP based semiconductor, and a current constriction layer 55b for the red element made of n-type GaAs which is the same material as that of the current constriction layer 55a, which is formed on sides of the ridge portion, and a contact layer 56b made of, for example, p-type GaAs is provided thereon. Further, p-side electrodes 57a and 57b are formed on the contact layers 56a and 56b, and a common n-side electrode 58 is formed at a back surface side of the semiconductor substrate 51, thereby a monolithic semiconductor laser is formed in which both elements are electrically separated. PATENT DOCUMENT 1: Japanese Patent Application Laid-Open No. 2000-11417 (FIG. 9)"}
{"text": "The present invention relates to a system for metering and dispensing, at a low flow rate, a high viscosity fluid, such as oil, and a method of metering and dispensing that high viscosity fluid.\nMotor oil that is used in automobiles, trucks, boats and other machinery has a very high viscosity. As used herein, the term \"high viscosity\" refers to fluids having a viscosity of in the range of 100 to 3,000 s.s.u., that is second say-bolt units. Currently, persons seeking to purchase quantities of oil would buy the oil in small units such as quarts or possibly single gallons. The oil is customarily sold in cans or quart containers (or gallon containers). This method of selling oil differs considerably from the sale of other petroleum products such as gasoline. In contrast, gasoline has a very low viscosity (approximately 30-40 s.s.u.) as compared with oil. As a reference, water has a viscosity of 30 s.s.u. Because of the compression problem when the oil includes air and because of the slippage problem when the oil flows through mechanical metering device, an oil dispensing system that meets government standards to controllably deliver set quantities of oil has not been commercially developed prior to the present invention.\nOne prior art device utilized a suction pump and an air eliminator but did not include an electronic calibration circuit that corrects for imperfections or characteristics in the oil dispensing system. In practice, it is necessary to calibrate to legal units of measure all government approved fluid dispensing systems that dispense fluid as a saleable product to the public. For example, the sale of gasoline to the public involves calibrating the pump, meter and the various system items. Since gasoline is a low viscosity and high flow rate product, compressibility of the fluid is not critical nor is slippage of the fluid through the system lines a problem. In contrast, a high viscosity fluid, such as oil has a high degree of slippage. Therefore, it is very difficult to design and manufacture a calibratable system to ensure that a customer receives a full legal quart of oil when a quart is purchased. The present invention overcomes these difficulties."}
{"text": "2-[4-(polyhalo-2-hydroxy-2-propyl)anilino]-2-oxazolines are known as disclosed in copending application Ser. No. 668,385, filed Mar. 19, 1976, now U.S. Pat. No. 4,081,547 as useful agents for the treatment of hypertension. The art also discloses that substituted anilino-2-oxazolines have hypotensive activity, e.g., U.S. Pat. Nos. 3,453,284; 3,499,083; and 3,499,084; and Belgian Pat. Nos. 704,392; 704,393; and 704,396. However, these compounds lack the 4-one or 2-thiazolin-4-one functional group present in the compounds of this invention."}
{"text": "The present invention is directed to bipolar transistor logic gates, and more particularly to gates of this type that combine high speed operation with high packing density.\nLogic gates that are constructed utilizing Schottky transistor logic (STL) are noted for their high switching speeds, and find utility in a number of applications because of this characteristic. Basically, an STL gate comprises a Schottky transistor, i.e., a bipolar transistor with a Schottky diode clamp between its base and collector. The cathodes of a plurality of Schottky diodes are each connected to the collector of the Schottky transistor, and their anodes comprise the output terminals for the gate. In operation, the clamp between the base and collector of the Schottky transistor prevents this transistor from going into saturation. Consequently, the transistor is quick to respond to changes in the voltage levels at its base. Further information regarding the features and operating characteristics of an STL gate can be found in the article \"STL Technology\" by Ben J. Sloan, IEDM Technical Digest 1979, pp. 324-327.\nWhile STL gates are desirable from the standpoint of their high switching speed, one drawback that has heretofore been associated with them is the comparatively large area that is required for each gate, resulting in a low packing density on a chip. More particularly, in a conventional STL gate, each of the output Schottky diodes is formed on a separate island on the chip, which increases the area of the buried layer for the chip and adds to the isolation area on the chip The higher the fan-out for each gate, the greater the number of islands that must be provided per gate. For example, a conventional STL gate having a fan-out of four requires five islands, one for each of the four Schottky diodes and one for the transistor collector. The need to provide each of these separate areas on the chip limits the packing density that can be obtained."}
{"text": "The invention relates to electronic entry control systems, and more particularly to self-contained electronic directory and entry control systems having liquid crystal directory displays, used to control access to apartment buildings and the like.\nIn apartment buildings, commercial office buildings, and other multi-tenant complexes, it is common practice to limit access to common areas to those individuals living and/or working within the complex, as well as their authorized guests and visitors. Of course, this can be accomplished in a simplistic manner by merely locking all access doors and providing tenants with keys. However, this approach has proven to be inconvenient and cumbersome in the admission of guests and visitors, as the tenant must physically unlock the door each time someone is to be admitted. Furthermore, with such a system, it is often difficult to determine whether a person is present at a door and desires admission. In the business environment, this can lead to a substantial loss in productivity.\nOne past attempt to resolve this problem has focused on the use of a dedicated intercom system. In such a system, a visitor actuates a button on an intercom panel outside the complex, alerting the desired tenant to his presence. Optionally, voice communication capability can be provided so that the visitor can establish his identity. Once it has been determined that the visitor is authorized to be admitted, the tenant can either physically unlock the door himself or press a button on the interior intercom panel to actuate a door unlocking mechanism.\nThe intercom system has several disadvantages. First, significant dedicated wiring is required to support the intercom at the entry door and to provide communication with each tenant unit. Furthermore, if remote door actuation capability is to be provided, further rewiring can be necessary. Expensive intercom hardware is required to be installed within each tenant unit. Finally, if many tenants reside within the complex, providing an individual \"call\" button for each tenant can result in a prohibitively large intercom unit at the entry door.\nIn response to the disadvantages of entry systems using dedicated intercoms, telephone-based entry systems have been proposed. Significant rewiring is not necessary, other than to provide electrical power and telephone line access to the station at the entry door. When a visitor desires access, he will place a telephone call from the entry system to the desired tenant, who will either physically unlock the door or, by using the touch-tone keypad on the telephone, provide an authorization signal to the entry system, which will then actuate a door unlocking mechanism. Additional expensive hardware is not needed for each tenant unit, as substantially all tenants have access to telephone equipment and a telephone line. Finally, since each tenant is addressable through a unique telephone number, a separate call button for each tenant is not necessary; numeric codes can be used.\nWith the resolution of the foregoing recognized problems by the use of telephone-based entry systems, other disadvantages have become apparent. Where a typical intercom system might have a separate call button for each individual tenant, each button can be labeled, making it easy to locate any desired tenant. Where no individual call buttons exist, however, a separate directory of tenants should be provided. Where there is a large number of tenants, the directory can be physically quite large. Moreover, labor involved in the maintenance of such directories to reflect tenant turnover can be significant, particularly when there are several entry doors to the complex.\nAccordingly, the need for an easily updatable electronic directory has been recognized. Such a goal is readily accomplished by using digital computer technology. As is well known, a small digital computer, with memory, can be used to store tenant names, telephone numbers, access codes, dialing codes, and other information within the entry system, at the expense of little additional cost or size. The electronic storage of such information makes it possible to update directory information with little or no effort; the revised information can be communicated to the digital computer through a data port, or even remotely via telephone lines.\nThe use of an electronic display, however, introduces further difficulties. Various display technologies are known and available for use. For example, a cathode ray tube (\"CRT\") video display, as is typically used for desktop computers, can easily be adapted for use as a directory display. However, CRT displays are generally large, heavy, and fragile, making them ill-suited for in-wall directory use. Furthermore, CRT displays consume a large amount of power. Light emitting diode (\"LED\") displays are also known, but are expensive and bulky when produced in sizes sufficient to display directory information. Fluorescent and gas-discharge displays, since they make use of a gas-filled glass tube, are also relatively fragile. Moreover, all of the foregoing display devices are light-emitting, which makes them difficult to see under bright light conditions, as in direct sunlight.\nLiquid crystal displays (\"LCDs\") are also known. They have the advantages of low cost and low power requirements. Such displays are selectively reflecting rather than light-emitting, so they are best viewed with ambient light, including bright light conditions that would \"wash out\" other types of displays. However, LCDs are also easily viewable in low light conditions when backlit.\nUnfortunately, LCDs are known to be temperature-sensitive. As temperature increases, the contrast of the display will often decline until it is unreadable. This is a particular problem where entry systems utilizing LCD directories are installed outdoors, in direct sunlight, or elsewhere under adverse temperature conditions. Furthermore, when an LCD is installed in a sealed enclosure with a viewing window, the temperature inside the enclosure may increase to a level far beyond the air temperature outside the enclosure due to a \"greenhouse effect.\"\nA variety of attempts have been made to solve the temperature sensitivity problem; none has proven to be entirely successful. One possibility is to provide a \"hood\" over the display to ensure that the display is shaded. This approach has the disadvantage of reducing the amount of light reaching the display, and it also limits the angles from which the display may be viewed. Furthermore, such an approach does not provide temperature reduction, per se; it merely reduces the possibility that direct sunlight from certain angles will raise the temperature.\nA second possibility is to ventilate the enclosure containing the LCD. However, this raises an opportunity for potential vandalism and damage to the entry system, as access to the interior of the enclosure is provided. Furthermore, ventilation of air to the interior of the enclosure will also allow the elements to invade the enclosure to some extent. Damage from oxidation, moisture, and other contaminants can accelerate the deterioration of the entry system and impair reliability. A further possibility is to provide a cooling fan. This, too, is unsatisfactory: moving parts are added to a system that could otherwise be entirely without them; furthermore, mechanical fans are noisy and potentially unreliable. When a fan is provided within a sealed enclosure, very little cooling action is possible.\nIt is recognized herein that the addition of an electronic display to an entry system complicates the system significantly. Adding the display unit will increase the size and cost of the entry system. Furthermore, if the display unit is of a size insufficient to display all tenant names at one time, a mechanism must be provided to determine which names are shown on the display. This can be accomplished by providing control buttons on the entry system, beyond those necessary for a simple telephonic keypad. Adding these additional controls will further raise the cost and complexity of the system, as well as the operational complexity to a typical user.\nSubstantially all practical systems of the type described herein providing communications capability between a visitor and a tenant, do so with audio communication. This type of communication is easily accomplished over the telephone system, and requires little additional circuitry or complexity to the system, beyond that already required to access the telephone system. Individuals who are deaf or hearing-impaired are not generally able to use these types of entry systems effectively.\nAccordingly, a need exists for an entry control system that avoids the many shortcomings of prior devices. Specifically, a need exists for a self-contained entry system that has a directory display that is easily visible under substantially all lighting conditions, and is resistant to adverse temperature conditions. A further need exists for an entry system that provides accessibility to deaf and hearing-impaired individuals. To reduce costs, such a system should minimize unneeded complexity."}
{"text": "The invention relates generally to remedial systems and methods for subsea structures. More particularly, the invention relates to systems and methods for pulling subsea structures such as primary conductors that have been bent from vertical.\nIn offshore drilling operations, subsea wells are built up by installing a primary conductor in the seabed and then securing a wellhead to the upper end of the primary conductor at the sea floor. A blowout preventer (BOP) is then installed on the wellhead, and a lower marine riser package (LMRP) mounted to the BOP. The primary conductor is typically installed in a vertical orientation to facilitate and simplify the installation of the BOP and LMRP onto the wellhead, which is coaxially aligned with the primary conductor. A lower end of a drilling riser is coupled to a flex joint on the top of the LMRP and extends to a drilling vessel or rig at the sea surface. A drill string is then suspended from the rig through the drilling riser, LMRP, BOP, wellhead, and primary conductor to drill a borehole while successively installing concentric casing strings that line the borehole. The casing strings are typically cemented at their lower ends and sealed with mechanical seals at their upper ends.\nDuring drilling operations, drilling fluid, or mud, is delivered through the drill string, and returned up an annulus between the drill string and casing that lines the borehole. In the event of a rapid influx of formation fluid into the annulus, commonly known as a “kick,” the BOP and/or LMRP may actuate to seal the annulus and control the well. In particular, BOPs and LMRPs comprise closure members capable of sealing and closing the well in order to prevent the release of high-pressure gas or liquids from the well. Thus, the BOP and LMRP are used as safety devices that close, isolate, and seal the wellbore. Heavier drilling mud may be delivered through the drill string, forcing fluid from the annulus through the choke line or kill line to protect the well equipment disposed above the BOP and LMRP from the high pressures associated with the formation fluid. Assuming the structural integrity of the well has not been compromised, drilling operations may resume. However, if drilling operations cannot be resumed, cement or heavier drilling mud is delivered into the well bore to kill the well.\nIn the event that the BOP and LMRP fail to actuate or insufficiently actuate in response to a surge of formation fluid pressure in the annulus, a blowout may occur. The blowout may damage subsea well equipment and hardware such as the BOP, LMRP, or drilling riser. For example, falling debris (e.g., a severed riser) resulting from a blowout may bend the primary conductor from the “as installed” vertical orientation. Bending of the primary conductor can also arise if the surface vessel drifts too far and exerts sufficiently large lateral loads on the LMRP and BOP via excessive tension applied to the riser extending from the surface vessel to the LMRP. In general, if the bending loads and associated stresses do not exceed the yield strength of the material forming the primary conductor, the primary conductor will not plastically deform and should rebound to its vertical orientation when the bending loads decrease sufficiently. However, if the bending loads and associated stresses exceed the yield strength of the material forming the primary conductor, the primary conductor will plastically deform and become permanently bent (i.e., the primary conductor will not rebound to its vertical orientation when the bending loads decrease)."}
{"text": "Gaming machines, such as slot machines, video poker machines and the like, have been a cornerstone of the gaming industry for several years. Generally, the popularity of such machines with players is dependent on the likelihood (or perceived likelihood) of winning money at the machine and the intrinsic entertainment value of the machine relative to other available gaming options. Where the available gaming options include a number of competing machines and the expectation of winning at each machine is roughly the same (or believed to be the same), players are likely to be attracted to the most entertaining and exciting machines Shrewd operators consequently strive to employ the most entertaining and exciting machines, features, and enhancements available because such machines attract frequent play and hence increase profitability to the operator. Therefore, there is a continuing need for gaming machine manufacturers to continuously develop new games and improved gaming enhancements that will attract frequent play through enhanced entertainment value to the player.\nOne concept that has been successfully employed to enhance the entertainment value of a game is the concept of a “secondary” or “bonus” game that may be played in conjunction with a “basic” game. The bonus game may comprise any type of game, either similar to or completely different from the basic game, which is entered upon the occurrence of a selected event or outcome in the basic game. Generally, bonus games provide a greater expectation of winning than the basic game and may also be accompanied with more attractive or unusual video displays and/or audio. Bonus games may additionally award players with “progressive jackpot” awards that are funded, at least in part, by a percentage of coin-in from the gaming machine or a plurality of participating gaming machines Because the bonus game concept offers tremendous advantages in player appeal and excitement relative to other known games, and because such games are attractive to both players and operators, there is a continuing need to develop gaming machines with new types of bonus games to satisfy the demands of players and operators.\nTraditional wagering games incentivize a player to maximize his or her own winnings, and emphasis is placed on appealing to the player's sense of self-gain. Even in tournament-style wagering games, each tournament player seeks to maximize his or her own winnings to the detriment of the other tournament players. Aspects disclosed herein address this and other related needs."}
{"text": "Spinal pathologies and disorders such as scoliosis and other curvature abnormalities, kyphosis, degenerative disc disease, disc herniation, osteoporosis, spondylolisthesis, stenosis, tumor, and fracture may result from factors including trauma, disease and degenerative conditions caused by injury and aging. Spinal disorders typically result in symptoms including deformity, pain, nerve damage, and partial or complete loss of mobility.\nNon-surgical treatments, such as medication, rehabilitation and exercise can be effective, however, may fail to relieve the symptoms associated with these disorders. Surgical treatment of these spinal disorders includes correction, fusion, fixation, discectomy, laminectomy and implantable prosthetics. As part of these surgical treatments, spinal constructs such as vertebral rods are often used to provide stability to a treated region. Rods redirect stresses away from a damaged or defective region while healing takes place to restore proper alignment and generally support the vertebral members. During surgical treatment, one or more rods may require bending for disposal with the vertebral members. Such rods may be attached via fasteners to the exterior of two or more vertebral members. This disclosure describes an improvement over these prior technologies."}
{"text": "A robot which transports a substrate such as a liquid crystal panel generally has a cantilevered hand with a fork on which a substrate is placed. In recent years, as the substrate size increases, the fork inevitably elongates. The fork itself is made of lightweight rigid carbon so as to reduce its flexure. In contrast, flexure occurs at the hand front end due to its own weight and the substrate weight to a larger degree in a direction away from the robot main body, that is, the center of gravity in the hand. Therefore, in the process of transporting a substrate to a processing apparatus for the next process or to a cassette, the substrate may cause a fault upon interfering with the processing apparatus or the cassette or may be damaged. To prevent this, various techniques for correcting flexure of the hand have been proposed.\nFor example, Japanese Patent Laid-Open No. 7-99225 discloses a technique for rotating an eccentric cam using an air cylinder to pivot a hand about a supporting portion in placing a substrate on the hand, thereby correcting the tilt angle of the hand in a direction opposite to that of flexure which has occurred upon placing the substrate on the hand.\nAlso, International Publication No. 2005/004230 discloses a technique for rotating a screw to tilt a tilting table on which a hand is fixed in placing a substrate on the hand, thereby correcting the tilt angle of the hand.\nNote that the amount of flexure of the hand changes depending on, for example, the weight of the transported substrate. In this respect, the technique disclosed in Japanese Patent Laid-Open No. 7-99225 assumes that the amount of flexure of the hand is constant. Therefore, this technique is inapplicable when substrates with different weights are transported. On the other hand, International Publication No. 2005/004230 discloses a method of storing in advance information concerning the amount of flexure of the hand due to the weight of each substrate, and correcting the tilt angle of the hand based on the type of transported substrate and the stored information. Unfortunately, it is cumbersome to store in advance information concerning the amount of flexure of the hand due to the weight of each substrate. Furthermore, it is impossible to correct the tilt angle of the hand in transporting a substrate whose information is not stored."}
{"text": "1. Field of the Invention\nThe present invention relates to a powdery pharmaceutical composition suitable for application to the mucosa of the oral or nasal cavity. More specifically, it relates to a powdery pharmaceutical composition, suitable for application to the mucosa of the oral or nasal cavity, containing (a) lower alkyl ethers of cellulose as a base, (b) a steroid or glycyrrhizic acid type anti-inflammatory agent as a drug component, and (c) a solid low irritant organic acid as a stabilizer for the drug component.\nThe powdery pharmaceutical composition exhibits extremely excellent stability in the steroid or glycyrrhizic acid type anti-inflammatory agent contained therein, low irritation of the mucosa of the oral or nasal cavity, and sustained release of the drug over an extended period of time. Accordingly, the powdery pharmaceutical composition according to the present invention can be effectively used for the treatment of oral diseases such as stomatitis, erosion or sores generated by radiotherapy, and licken planus or nasal diseases such as allergic rhinitis and vasomotor rhinitis.\n2. Description of the Related Art\nVarious pharmaceutical preparations, such as buccal tablets, troche tablets, sublingual tablets, and oral ointments, have been heretofore known in the art for application to the oral cavity. Buccal tablets, troche tablets, and sublingual tablets, however, give a foreign feeling to the oral cavity. Therefore, patients often chew on and swallow the tablets. Oral ointments, on the other hand, usually contain mixtures of beeswax or plastibase with gelatin, pectin, or sodium carboxymethyl cellulose. These mixtures are unsatisfactory because the adhesion properties thereof to the oral mucosa are not sufficient and also because they have unpleasant tastes.\nOn the other hand, various pharmaceutical preparations, such as nasal ointments, jellies, nose drops, and sprays, have been known in the art for application to the nasal cavity. Nasal ointments and jellies, however, are difficult to apply to deep parts of the nasal cavity, such as the concha nasalis superior. With, nose drops and sprays, the active components do not remain in the nasal cavity for an extended period of time.\nVarious sustained release pharmaceutical preparations for oral or nasal cavity administration are known in the art. With these, the efficacy of the drugs is spread over a long period of time by the gradual or sustained releasing the active drug components from the pharmaceutical preparations. For example, Japanese Unexamined Patent Publication (Kokai) No. 57-118511 discloses powdery sustained release preparations suitable for adhesion to the oral cavity. These preparations use, as a base, hydroxypropyl cellulose and have the advantages of an excellent sustained release of the drugs and possible application over a wide range of the oral cavity. Furthermore, U.S. Pat. No. 4,294,829 discloses powdery sustained release preparations suitable for application to the nasal cavity. These contain lower alkyl ethers of cellulose and drugs and also have the advantage of excellent sustained release of the drugs.\nIt is, however, more desirable to improve the stability of drugs in the above-mentioned sustained release preparations for application to the mucosa of the oral or nasal cavity."}
{"text": "The invention relates to methods and apparatus for the provision of ventilatory assistance matched to a subject\"\"s respiratory need. The ventilatory assistance can be for a subject who is either spontaneously or non-spontaneously breathing, or moves between these breathing states. The invention is especially suitable for, but not limited to, spontaneously breathing human subjects requiring longterm ventilatory assistance, particularly during sleep\nSubjects with severe lung disease, chest wall disease, neuromuscular disease, or diseases of respiratory control may require in-hospital mechanical ventilatory assistance, followed by longterm home mechanical ventilatory assistance, particularly during sleep. The ventilator delivers air or air enriched with oxygen to the subject, via an interface such as a nosemask, at a pressure that is higher during inspiration and lower during expiration.\nIn the awake state, and while waiting to go to sleep, the subject\"\"s ventilatory pattern is variable in rate and depth. Most known ventilatory devices do not accurately match the amplitude and phase of mask pressure to the subject\"\"s spontaneous efforts, leading to discomfort or pac. Larger amounts of asynchrony also reduce the efficiency of the device. During sleep, thee are changes in the neural control of breathing as well as the mechanics of the subject\"\"s airways, respiratory muscles and chest wall, leading to a need for substantially increased ventilatory support. Therefore, unless the device can automatically adjust the degree of support, the amplitude of delivered pressure will either be inadequate during sleep, or must be excessive in the awake state. This is particularly important in subjects with abnormalities of respiratory control, for example central hypoventilation syndromes, such as Obesity Hypoventilation Syndrome, where there is inadequate chemoreceptor drive, or Cheyne Stokes breathing such as in patients with severe cardiac failure or after a stroke, where there is excessive or unstable chemoreceptor drive.\nFurthermore, during sleep there are inevitably large leaks between mask and subject, or at the subject\"\"s mouth if this is left free. Such leaks worsen the error in matching the phase and magnitude of the machine\"\"s effort to the subject\"\"s needs, and, in the case of mouth leak, reduce the effectiveness of the ventilatory support.\nIdeally a ventilatory assistance device should simultaneously address the following goals:\n(i) While the subject is awake and making substantial ventilatory efforts, the delivered assistance should be closely matched in phase with the patient\"\"s efforts.\n(ii) The machine should automatically adjust the degree of assistance to maintain at least a specified ventilation, without relying on the integrity of the subject\"\"s chemoreflexes.\n(iii) It should continue to work correctly in the presence of large leaks.\nMost simple home ventilators either deliver a fixed volume, or cycle between two fixed pressures. They do so either at a fixed rate, or are triggered by the patient\"\"s spontaneous efforts, or both. All such simple devices fail to meet goal (ii) of adjusting the degree of assistance to maintain at least a given ventilation. They also largely fail to meet goal (i) of closely matching the subjects respiratory phase: timed devices make no attempt to synchronize with the subject\"\"s efforts; triggered devices attempt to synchronize the start and end of the breath with the subject\"\"s efforts, but make no attempt to tailor the instantaneous pressure during a breath to the subject\"\"s efforts. Furthermore, the triggering tends to fail in the presence of leaks, thus failing goal (iii).\nThe broad family of servo-ventilators known for at least 20 years measure ventilation and adjust the degree of assistance to maintain ventilation at or above a specified level, thus goal (ii), but they still fail to meet goal (i) of closely matching the phase of the subject\"\"s spontaneous efforts, or the reasons given above. No attempt is made to meet goal (iii).\nProportional assistist ventilation (PAV), as taught by Dr Magdy Younes, for example in Principles wad Practice of Mechanical Ventilation, chapter 15, aims to tailor the pressure vs time profile within a breath to partially or completely unload the subject\"\"s resistive and elastic work, while minimizing the airway pressure required to achieve the desired ventilation. During the inspiratory half-cycle, the administred pressure takes the form:\nP(t)=P0+R.fRESP(t)+E.V(t)\nwhere R is a percentage of the resistance of the airway, fRESP(t) is the instantaneous respiratory airflow at time t, E is a percentage of the elastance of lung and chest wall, and V(t) is the volume inspired since the start of inspiration to the present moment. During the expiratory half-cycle, V(t) is taken as zero, to produce passive expiration.\nAn advantage of proportional assist ventilation during spontaneous breathing is that the degree of assistance is automatically adjusted to suit the subject\"\"s immediate needs and their pattern of breathing, and is therefore comfortable in the spontaneously breathing subject. However, there are at least two important disadvantages. Firstly, V(t) is calculated as the integral of flow with respect to time since the start of inspiration. A disadvantage of calculating V(t) in this way is that, in the presence of leaks, the integral of the flow through the leak will be included in V(t), resulting in an overestimation of V(t), in turn resulting in a runaway increase in the administered pressure. This can be distressing to the subject. Secondly, PAV relies on the subject\"\"s chemoreceptor reflexes to monitor the composition of the arterial blood, and thereby set the level of spontaneous effort, The PAV device then amplifies this spontaneous effort. In subjects with abnormal chemoreceptor reflexes, the spontaneous efforts may either cease entirely, or become unrelated to the composition of the arterial blood, and amplification of these efforts will yield inadequate ventilation. In patients with existing Cheyne Stokes breathing during sleep, PAV will by design amplify the subject\"\"s waxing and waning breathing efforts, and actually make matters worse by exaggerating the disturbance. Thus PAV substantially meets goal (i) of providing assistance in phase with the subject\"\"s spontaneous ventilation, but cannot meet goal (ii) of adjusting the depth of assistance if the subject has inadequate chemoreflexes, and does not satisfactorily meet goal (iii).\nThus there are known devices that meet each of the above goals, but there is no device hat meets all the goals simultaneously. Additionally, it is desirable to provide improvements over the prior art directed to any one of the stated goals.\nTherefore, the present invention seeks to achieve, at least partially, one or more of the following:\n(i) to match the phase and degree of assistance to the subject\"\"s spontaneous efforts when ventilation is well above a target ventilation,\n(ii) to automatically adjust the degree of assistance to maintain at least a specified minimum average ventilation without relying on the integrity of the subject\"\"s chemoreflexes and to damp out instabilities in the spontaneous ventilatory efforts, such as Cheyne Stokes breathing.\n(iii) to provide some immunity to the effects of sudden leaks.\nIn what follows, a fuzzy membership function is taken as returning a value between zero and unity, fit intersection A AND B is the smaller of A and B, fuzzy union A OR B is the larger of A and B, and fuzzy negation NOT A is 1-A.\nThe invention discloses the determination of the instantaneous phase in the respiratory cycle as a continuous variable.\nThe invention further discloses a method for calculating the instantaneous phase in dew respiratory cycle including at least the steps of determining that if the instantaneous airflow is small and increasing fast, then it is close to start of inspiration, if the instantaneous airflow is large and steady, then it is close to mid-inspiration, if the instantaneous airflow is small and decreasing fast, then it is close to mid-expiration, if the instantaneous airflow is zero and steady, then it is during an end-expiratory pause, and airflow conditions intermediate between the above are associated with correspondingly intermediate phases.\nThe invention further discloses a method for determining the instantaneous phase in the respiratory cycle as a continuous variable from 0 to 1 revolution, the method comprising the steps of:\nselecting at least two identifiable features FN of a prototype flow-vs-time waveform ƒ(t) similar to an expected respiratory flow-vs-time waveform, and for each said feature:\ndetermining by inspection the phase xcfx86N in the respiratory cycle for said feature, assigning a weight WN to said phase,\ndefining a xe2x80x9cmagnitudexe2x80x9d fuzzy set MN whose membership function is a function of respiratory airflow, and a xe2x80x9crate of changexe2x80x9d fuzzy set CN, whose membership function is a function of the time derivative of respiratory airflow, chosen such that the fuzzy intersection MN AND CN will be larger for points on the generalized prototype respiratory waveform whose phase is closer to the said feature FN than for points closer to all other selected features,\nsetting the fuzzy inference rule RN for the selected feature FN to be: If flow is MN and rate of change of flow is CN then phase=xcfx86N, with weight WN,\nmeasuring leak-corrected respiratory airflow,\nfor each feature FN calculating fuzzy membership in fuzzy sets MN and CN,\nfor each feature FN applying fuzzy inference rule RN to determine the fuzzy extent YN=MN AND CN to which the phase is xcfx86N, and\napplying a defuzzification procedure using YN at phases xcfx86N and weights WN to determine the instantaneous phase xcfx86.\nPreferably, the identifiable features include zero crossings, peaks, inflection points or plateaus of the prototype flow-vs-time waveform. Furthermore, said weights can be unity, or chosen to reflect the anticipated reliability of deduction of the particular feature.\nThe invention further discloses a method for calculating instantaneous phase in the respiratory cycle as a continuous variable, as described above, in which the step of calculating respiratory airflow includes a low pass filtering step to reduce non-respiratory noise, in which the time constant of the low pass filter is an increasing function of an estimate of the length of the respiratory cycle.\nThe invention further discloses a method for measuring the instantaneous phase in the respiratory cycle as a continuous variable as described above, in which the defuzzification step includes a correction for any phase delay introduced in the step of low pass filtering respiratory airflow.\nThe invention further discloses a method for measuring the average respiratory rate, comprising the steps of:\nmeasuring leak-corrected respiratory airflow,\nfrom the respiratory airflow, calculating the instantaneous phase xcfx86 in the respiratory cycle as a continuous variable from 0 to 1 revolution, calculating the instantaneous rate of change of phase dxcfx86/dt, and\ncalculating the average respiratory rate by low pass filtering said instantaneous rate of change of phase dxcfx86/dt.\nPreferably, the instantaneous phase is calculated by the methods described above.\nThe invention further discloses a method for providing ventilatory assistance in a spontaneously breathing subject, comprising the steps, performed at repeated sampling intervals, of:\nascribing a desired waveform template function Π(xcfx86), with domain 0 to 1 revolution and range 0 to 1,\ncalculating the instantaneous phase xcfx86 in the respiratory cycle as a continuous variable from 0 to 1 revolution,\nselecting a desired pressure modulation amplitude A,\ncalculating a desired instantaneous delivery pressure as an end expiratory pressure plus the desired pressure modulation amplitude A multiplied by the value of the waveform template function Π(xcfx86) at the said calculated phase xcfx86, and\nsetting delivered pressure to subject to the desired delivery pressure.\nThe invention further discloses a method for providing ventilatory assistance in a spontaneously breathing subject as described above, in which the step of selecting a desired pressure modulation amplitude is a fixed amplitude.\nThe invention further discloses a method for providing ventilatory assistance in a spontaneously breathing subject as described above, in which the step of selecting a desired pressure modulation amplitude in which said amplitude is equal to an elastance multiplied by an estimate of the subject\"\"s tidal volume.\nThe invention further discloses a method for providing ventilatory assistance in a spontaneously breathing subject as described above, in which the step of selecting a desired pressure modulation amplitude comprises the substeps of:\nspecifying a typical respiratory rate giving a typical cycle time,\nspecifying a preset pressure modulation amplitude to apply at said typical respiratory rate,\ncalculating the observed respiratory rate giving an observed cycle time, and\ncalculating the desired amplitude of pressure modulation as said preset pressure modulation amplitude multiplied by said observed cycle time divided by the said specified cycle time.\nThe invention further discloses a method for providing ventilatory assistance in a spontaneously breathing subject, including at least the step of determining the extent that the subject is adequately ventilated, to said extent the phase in the respiratory cycle is determined from the subject\"\"s respiratory airflow, but to the extent that the subject\"\"s ventilation is inadequate, the phase in the respiratory cycle is assumed to increase at a pre-set rate, and setting mask pressure as a function of said phase.\nThe invention further discloses a method for providing ventilatory assistance in a spontaneously breathing subject, comprising the steps of: measuring respiratory airflow, determining the extent to which the instantaneous phase in the respiratory cycle can be determined from said airflow, to said extent determining said phase from said airflow but to the extent that the phase in the respiratory cycle cannot be accurately determined, the phase is assumed to increase at a preset rate, and delivering pressure as a function of said phase.\nThe invention further discloses a method for calculating the instantaneous inspired volume of a subject, operable substantially without run-away under conditions of suddenly changing leak, the method comprising the steps of:\ndetermining respiratory airflow approximately corrected for leak,\ncalculating an index J varying from 0 to 1 equal to the fuzzy extent to which said corrected respiratory airflow is large positive for longer than expected, or large-negative for longer than expected,\nidentifying the start of inspiration, and\ncalculating the instantaneous inspired volume as the integral of said corrected respiratory airflow multiplied by the fuzzy negation of said index J with respect to time, from start of inspiration.\nThe invention further discloses a method xe2x80x9cAxe2x80x9d for providing ventilatory assistance in a spontaneously breathing subject, the method comprising the steps, performed at repeated sampling intervals, of:\ndetermining respiratory airflow approximately corrected for leak,\ncalculating an index J varying from 0 to 1 equal to the fuzzy extent to which said respiratory airflow is large positive for longer than expected, or large negative for longer than expected,\ncalculating a modified airflow equal to said respiratory airflow multiplied by the fuzzy negation of said index J.\nidentifying the phase in the respiratory cycle,\ncalculating the instantaneous inspired volume as the integral of said modified airflow with respect to time, with the integral held at zero during the expiratory portion of the respiratory cycle,\ncalculating a desired instantaneous delivery pressure as a function at least of the said instantaneous inspired volume, and\nsetting delivered pressure to subject to the desired delivery pressure.\nThe invention further discloses a method xe2x80x9cBxe2x80x9d for providing ventilatory assistance in a spontaneously breathing subject, comprising the steps of:\ndetermining respiratory airflow approximately corrected for leak,\ncalculating an index J varying from 0 to 1 equal to the fuzzy extent to which the respiratory airflow is large positive for longer than expected, or large negative for longer than expected,\nidentifying the phase in the respiratory cycle,\ncalculating a modified respiratory airflow equal to the respiratory airflow multiplied by the fuzzy negation of said index J,\ncalculating the instantaneous inspired volume as the integral of the modified airflow with respect to time, with the integral held at zero during the expiratory portion of the respiratory cycle,\ncalculating the desired instantaneous delivery pressure as an expiratory pressure plus a resistance multiplied by the instantaneous respiratory airflow plus a nonlinear resistance multiplied by the respiratory airflow multiplied by the absolute value of the respiratory airflow plus an elastance multiplied by the said adjusted instantaneous inspired volume, and\nsetting delivered pressure to subject to the desired delivery pressure.\nThe invention yet further discloses a method xe2x80x9cCxe2x80x9d for providing assisted ventilation to match the subject\"\"s need, comprising the steps of:\ndescribing a desired waveform template function Π(xcfx86), with domain 0 to 1 revolution and range 0 to 1,\ndetermining respiratory airflow approximately corrected for leak,\ncalculating an index J varying from 0 to 1 equal to the fuzzy extent to which the respiratory airflow is large positive for longer than expected, or large negative for longer than expected,\ncalculating JPEAK equal to the recent peak of the index J,\ncalculating the instantaneous phase in the respiratory cycle,\ncalculating a desired amplitude of pressure modulation, chosen to servo-control the degree of ventilation to at least exceed a specified ventilation,\ncalculating a desired delivery pressure as an end expiratory pressure plus the calculated pressure modulation amplitude A multiplied by the value of the waveform template function Π(xcfx86) at the said calculated phase xcfx86, and\nsetting delivered pressure to subject to said desired instantaneous delivered pressure,\nThe invention yet further discloses a method for providing assisted ventilation to match the subject\"\"s need, as described above, in which the step of calculating a desired amplitude of pressure modulation, chosen to servo-control the degree of ventilation to at least exceed a specified ventilation, comprises the steps of:\ncalculating a target airflow equal to twice the target ventilation divided by the target respiratory rate,\nderiving an error term equal to the absolute value of the instantaneous low pass filtered respiratory airflow minus the target airflow, and\ncalculating the amplitude of pressure modulation as the integral of the error term multiplied by a gain, with the integral clipped to lie between zero and a maximum.\nThe invention yet further discloses a method for providing assisted ventilation to match the subject\"\"s need, as described above, in which the step of calculating a desired amplitude of pressure modulation, chosen to servo-control the degree of ventilation to at least exceed a specified ventilation, comprises the following steps:\ncalculating a target airflow equal to twice the target ventilation divided by the target respiratory rate,\nderiving an error term equal to the absolute value of the instantaneous low pass filtered respiratory airflow minus the target airflow,\ncalculating an uncorrected amplitude of pressure modulation as the integral of the error term multiplied by a gain, with the integral clipped to lie between zero and a maximum,\ncalculating the recent average of said amplitude as the low pass filtered amplitude, with a time constant of several times the length of a respiratory cycle, and\nsetting the actual amplitude of pressure modulation to equal the said low pass filtered amplitude multiplied by the recent peak jamming index JPEAK plus the uncorrected amplitude multiplied by the fuzzy negation of JPEAK.\nThe invention yet further discloses a method for providing assisted ventilation to match the subject\"\"s need, and with particular application to subjects with varying respiratory mechanics, insufficient respiratory drive, abnormal chemoreceptor reflexes, hypoventilation syndromes, or Cheyne Stokes breathing, combined with the advantages of proportional assist ventilation adjusted for sudden changes in leak, comprising the steps, performed at repeated sampling intervals, of:\ncalculating the instantaneous mask pressure as described for methods xe2x80x9cAxe2x80x9d or xe2x80x9cBxe2x80x9d above,\ncalculating the instantaneous mask pressure as described for method xe2x80x9cCxe2x80x9d above,\ncalculating a weighted average of the above two pressures, and\nsetting the mask pressure to the said weighted average.\nThe invention yet further discloses apparatus to give effect to each one of the methods defined, including one or more transducers to measure flow and/or pressure, processor means to perform calculations and procedures, flow generators for the supply of breathable gas at a pressure above atmospheric pressure and gas delivery means to deliver the breathable gas to a subject\"\"s airways.\nThe apparatus can include ventilators, ventilatory assist devices, and CPAP devices including constant level, bi-level or autosetting level devices.\nIt is to be understood that while the algorithms embodying the invention are explained in terms of fuzzy logic, approximations to these algorithms can be constructed without the use of the fuzzy logic formalism."}
{"text": "1. Field of the Invention\nThis invention generally relates to cable networks and to telecommunications networks and, more particularly, to methods and systems that receive electrical signals from a cable network and communicate those electrical signals along a telecommunications network.\n2. Description of the Related Art\nTelecommunications is converging into broadband networks. Telecommunications networks in the past have been generally divided into three distinct groups: i) community antenna television networks, ii) local and wide-area computer networks, and iii) public-switched telephone networks. These three networks were independent of each other and delivered unique services to customers. The traditional cable networks, for example, delivered one-way video content, however, the cable networks did not offer telephony and data services. Telephone networks, similarly, did not offer analog and digital video services. Today, however, the Telecommunications Act of 1996 ushered in deregulation and permitted any telecommunications service provider to enter another's market. Because of this legislative enactment, the three separate telecommunications services are converging into a single broadband network and, thus, delivering a wide range of entertainment and data services.\nWhile this convergence benefits consumers, this convergence presents problems for the telecommunications service provider. If a traditional telephone service provider, for example, wants to provide digital video and music content to their subscribers, the telephone service provider must first obtain the digital content. That is, the telephone service provider must obtain the music, videos, and other content to provide to customers. Obtaining this content, however, requires a heavy investment in an entertainment infrastructure. The telephone service provider must find and employ people with knowledge and experience in the entertainment industry. The telephone service provider must determine who is the target subscriber and what content satisfies that target subscriber. The telephone service provider must also obtain the intellectual property rights to legally provide this content. These simple but broad concerns often represent a very expensive investment in an infrastructure that provides content along a converged telecommunications network. There is, accordingly, a need for a less expensive alternative to converged networks, a need for telecommunications service providers to less expensively provide broadband content, and, in particular, a need for cable networks and telephone networks to provide broadband content without a large investment in separate and competing networks."}
{"text": "1. Field of the Invention\nThis invention relates to proteins which inhibit the coagulation of the blood, processes for preparing these proteins, and their use.\n2. Description of the Background Art\nAnti-coagulant proteins, which are present in most mammals, can be divided into three groups based on their different mechanisms of activity.\nOne group of proteins form a complex with a coagulation factor and thereby render the coagulation factor inactive. Proteins in this category include antithrombin III (Thromb. Res., 5:,439-452 (1974)), alpha.sub.1 -protease inhibitor (Ann. Rev. Biochem., 52:655-709 (1983)), alpha.sub.2 -macroglobulin (Ann. Rev Biochem., 52: 655-709 (1983)), C.sub.1 -inhibitor (Biochemistry 20: 2738-2743 (1981)), and protease nexin (J. Biol. Chem., 258: 10439-10444 (1983)).\nA second group of proteins act proteolytically on a coagulating factor and thereby inactivate it. The only protein of this kind that has been described is protein C (J. Biol. Chem., 251: 355-363 (1976)).\nThe third category to which anti-coagulant proteins can be grouped are those which screen and/or hydrolyze the negatively charged phospholipids so that the phospholipid-dependent reactions of the blood coagulation mechanism are inhibited. Thus far, only phospholipases isolated from various types of snake venom have been described as having this mode of action (Eur. J. Biochem., 112: 25-32 (1980)).\nIn recent years, the step-wise coagulation system has been investigated thoroughly. It is understood to be an intensifying multi-stage system of different interconnected proteolytic reactions in which an enzyme converts a zymogen into the active form (cf. Jackson, C. M. and Nemerson, Y., Ann. Rev. Biochem., 49: 765-811 (1980)). The speed of this reaction is decisively increased by the presence of phospholipids and other cofactors such as factor V.sub.a and factor VIII.sub.a. In vivo, the procoagulation reactions are regulated by a variety of inhibitory mechanisms which prevent an explosively thrombotic trauma after slight activation of the coagulation cascade.\nThe mechanisms by which the anti-coagulation proteins of these three groups act have been described (Rosenberg, R. D. and Rosenberg, J. S., J. Clin. Invest., 74: 1-6 (1984)).\nIn Group 1, serine-protease factor X.sub.a and thrombin are inactivated as a result of their binding to antithrombin III or to the antithrombin/heparin complex. Both the prothrombin activation and also the formation of fibrin can be inhibited in this way. In addition to antithrombin III, there are also various other plasmaprotease inhibitors such as alpha2-macroglobulin and antitrypsin, the activity of which is dependent on time.\nIn Group 2, the discovery of protein C led to another anti-coagulation mechanism. Once protein C is activated, it acts as an anti-coagulant by selective proteolysis of the protein cofactors V.sub.a and VIII.sub.a, by which prothrombinase and the enzyme which converts factor X are deactivated.\nIn Group 3, plasmin cleaves monomeric fibrin 1, a product of the effect of thrombin on fibrinogen, thereby preventing the formation of an insoluble fibrin (Nossel, H. L., Nature, 291: 165-167 (1981)).\nOf the above-mentioned native proteins involved in the coagulation process, at present only antithrombin III is clinically used. However, the increase in the tendency to bleed when this protein is administered has proven to be a serious disadvantage.\nAll the agents previously used as anticoagulants, whether native to the body or synthetic, in some way render the coagulation factors ineffective and thereby lead to side effects which have a disadvantageous effect on the coagulation process."}
{"text": "Gene therapy by delivery of transgenes encoding therapeutic neurotrophic factors into the brain offers great promise for treating neurodegenerative conditions such as Alzheimer's Disease, Parkinson's Disease and Huntington's Disease. The protocols by which these therapies are provided are highly exacting, requiring that the therapeutic composition dosage (e.g., viral titer) be consistently provided at precise locations in the brain, to ensure that a predictable amount of neurotrophic factor be delivered only to targeted cells.\nFor example, U.S. Pat. No. 6,451,306 provides a method for treating Alzheimer's Disease which requires donor cells containing a neurotrophic factor-encoding transgene to be grafted at pre-determined sites in the forebrain located no more than 5 mm apart and no more than 500 μm from a targeted cell. The dosage of donor cells provided at each site preferably falls within a range of 2 to 20 μl per ml of composition. Similarly, U.S. Pat. Nos. 6,683,058 and 6,851,431 provide methods to treat defects, disease or damaged cholinergic and dopaminergic neuron populations, respectively, by delivering transgenes at sites within 500 μm of a targeted neuron and no more than 10 mm apart. Such parameters leave the practicing neurosurgeon relatively little room for error in dosing or placement of each graft or transgene injection.\nYet the conventional drug delivery devices available for use in gene therapy of the brain do not necessarily provide the consistent precision the therapeutic protocols require. For example, it has been reported that polynucleotides can become inactivated when introduced through a conventional needle cannula; e.g., made of a metal such as stainless steel (see, e.g., U.S. Pat. No. 7,060,056). The '056 patent inventors opined that the metal interacted with polynucleotides in a way that compromised their pharmaceutical activity by inactivation, rather than binding (see, e.g., '056 patent, Example 10). However, the inventors have discovered that many commonly used viral recombinant expression vectors are not inactivated by contact with metals and the like, most likely due to their proteinaceous coating (e.g., the capsid proteins of a virus). Instead, such material retains its activity, but is lost during delivery to binding within the lumen of conventional metal needles. The extent of loss varies from passage to passage, which limits the clinician's ability to accurately predict how much viral vector will actually be delivered out of any given injection or infusion.\nThe margin for dosing error in gene therapy of the brain can be increased if the instruments utilized to deliver a neurotrophic factor-encoding transgene (e.g., as part of a viral vector) cannot be consistently and accurately targeted to cell populations that may be only microns apart. If a target cell is missed, the extent to which expressed neurotrophic factor secreted by another cell will diffuse to a targeted cell is limited. Therefore, improvements in therapeutic efficacy can be obtained by enhancing the accurate placement of transgene-containing donor cell grafts or viral vectors into the brain."}
{"text": "Digital watermarking is a process for modifying physical or electronic media to embed a hidden machine-readable code into the media. The media may be modified such that the embedded code is imperceptible or nearly imperceptible to the user, yet may be detected through an automated detection process. Most commonly, digital watermarking is applied to media signals such as images, audio signals, and video signals. However, it may also be applied to other types of media objects, including documents (e.g., through line, word or character shifting), software, multi-dimensional graphics models, and surface textures of objects.\nDigital watermarking systems typically have two primary components: an encoder that embeds the watermark in a host media signal, and a decoder that detects and reads the embedded watermark from a signal suspected of containing a watermark (a suspect signal). The encoder embeds a watermark by subtly altering the host media signal. The reading component analyzes a suspect signal to detect whether a watermark is present. In applications where the watermark encodes information, the reader extracts this information from the detected watermark.\nSeveral particular watermarking techniques have been developed. The reader is presumed to be familiar with the literature in this field. Particular techniques for embedding and detecting imperceptible watermarks in media signals are detailed in the assignee's U.S. Pat. Nos. 6,122,403 and 6,614,914, which are hereby incorporated by reference.\nIntegrating a hardware watermark reader in imaging devices (cameras, scanners, printers, fax machines, copiers) cellular phones, PDA devices may soon be necessary. Currently, most of these devices have no digital rights management (DRM) and hence can freely download music, video, and pictures from the Internet. An integrated watermark reader not only would enforce DRM in these devices, but also it would enable new applications. Such applications may potentially change the usage habit and pattern of these devices. For example, the imaging device integrated in most of these devices could be used with a watermarked advertisement and a built-in global positioning system (GPS) to guide the user to the nearest store where the advertised goods are currently available or on sale.\nLow-cost stand-alone watermark readers will also be developed soon to harness the power of the millions of watermark-enhanced identity cards already deployed in the USA and around the world. These readers will be essential in enforcing the law and reducing identity thefts. Also, the ubiquitous deployment of these readers will enable and empower a new wave of applications.\nCellular phones and PDA devices, however, have some limitations when used with sophisticated watermark detectors operating on large images. These limitations include the type and speed of the processor used in these devices, the amount of available memory, the system bandwidth, and the operating system controlling these devices.\nUsually, cellular phones and PDA devices employ fixed-point processors such as Intel's StrongARM or XScale processors, which have a clock speed in the 100-400 MHz range. This speed is only a fraction of the speed of mainstream processors typically found in PCs. This speed limits the amount of computations and their precision, which the watermark detector is allowed to perform.\nThese devices also have limited DRAM and system bandwidth. Currently, PDA devices have 32 MB of memory, but cellular phones have about 6 MB only. About half of the memory available on cellular phones is used by the operating system. Also, most of these devices have about 20 MHz of system bandwidth. The low bandwidth combined with the lack of memory can be a major bottleneck for real-time detection of the watermarks embedded in large images via cellular phones and PDA devices.\nMoreover, these devices use either a primitive native operating system such as MS SmartPhone, Symbian, and Brew, or a Java-based operating system such as Doja or J2ME. Although native operating systems allow the execution of a watermark detector written in a compiled language such as C, Java-based operating systems require Java implementations of the detector. Such an implementation is not as efficient as a C implementation, especially for sophisticated signal processing algorithms such as watermark detectors.\nAlthough a watermark detector can be implemented in a customized chip, the cost of implementing a complex algorithm can be prohibitive. A field programmable gate array (FPGA) or application-specific integrated circuit (ASIC) can be used for this purpose. Such a chip can be integrated with the cellular phone, PDA device, or stand-alone watermark reader. However, complex watermarking algorithms may require a large number of gates and excessive on-chip memory and system bandwidth, which tremendously increases the cost of this dedicated chip. This cost, however, can be reduced substantially by reducing the data that the watermark detector operates on and the computations it performs on this data.\nThis document describes a block-based watermark detector and reader that use a hierarchical search strategy to quickly zoom into the image region with the strongest watermark. This approach reduces the computations and the amount of data the detector operates on, which, in turn, reduces the processor speed, the memory, and the bandwidth requirement of the detector. This design makes the detector suitable for the aforementioned real-time software and low-cost hardware implementations.\nFurther features and combinations will become even more apparent with reference to the following detailed description and accompanying drawings."}
{"text": "For an uplink channel of 3rd generation partnership project radio access network long term evolution (3GPP RAN LTE, hereinafter referred to as “LTE”), single-carrier transmission is adapted to achieve a low peak to average power ratio (PAPR).\nFurther, for an LTE uplink channel, to achieve high throughput, adaptive modulation (AMC: adaptive modulation and coding) is used so as to select the modulation and coding scheme (MCS) pattern for each user depending on the channel quality indicator (CQI) of each user.\nFurther, introduction of the MIMO system is being considered to achieve higher transmission rate and to further improve the efficiency of use of frequency. Introduction of a rank transmission technique is also considered such as rank adaptation with which the rank indication (the number of rank) is adaptively switched depending on the status of a spatial propagation path to further improve transmission rate.\nUnder these circumstances, an agreement is made to time-multiplex control information and user data using the physical uplink shared channel (PUSCH) of the same subframe so as to maintain low PAPR even when both control information and user data are transmitted at the same time in an LTE uplink channel (see Non-Patent Literature 1).\nThe number of coded symbols, Q′, of control information to be multiplexed with user data is set based on equation 1.\n                                              ⁢                  [          1          ]                                                                                              ⁢                              Q            ′                    =                                    min              (                                                ⌈                                      O                                                                  10                                                                              -                                                          Δ                              offset                              PUSCH                                                                                10                                                                    .                                              R                        data                                                                              ⌉                                ,                                  4                  ·                                      M                    sc                                                              )                        =                          min              ⁡                              (                                                      Q                    ⁢                                                                                  ⁢                    1                                    ,                                      Q                    ⁢                                                                                  ⁢                    2                                                  )                                                                        (                  Equation          ⁢                                          ⁢          1                )            where ┌x┐ is an integer not greater than x, and min(x,y) is the value of the smaller one of x and y.\nIn equation 1, Msc is the number of subcarriers per PUSCH subframe, and ΔPUSCHoffset is a PUSCH offset which varies per control information such as ACK/NACK, rank indicator (RI), or CQI. ΔPUSCHoffset is reported from a higher layer (see Non-Patent Literature 1).\nO is the number of control information bits, and Rdata is represented by equation 2.\n                    [        2        ]                                                                      R          data                =                                            ∑                              r                =                0                                            C                -                1                                      ⁢                          K              r                                                          M              sc                        ·                          N              symb                                                          (                  Equation          ⁢                                          ⁢          2                )            In equation 2, K is the number of bits on r-th block, C is the number of blocks per PUSCH subframe, and Nsymb is the number of symbols per PUSCH subcarrier. The actual coding rate of user data is obtained by dividing Rdata in equation 2 by the number of bits per symbol, and is in proportion to Rdata equation 2. Accordingly, Rdata in equation 2 will be hereinafter referred to as “user data coding rate.”\nIn equation 1, Q1 is the number of coded symbols of control information that is set based on the number of control information bits O, user data coding rate Rdata, and PUSCH offset ΔPUSCHoffset, per control information. Q2 is the upper limit value of the number of coded symbols of control information. As shown in equation 1, the number of coded symbols of control information, Q′, is set by the smaller one of the number of symbols Q1 and upper limit value Q2.\nHere, equation I is modified to provide equation 3. As is the case with Rdata in equation 2, the actual coding rate of control information is obtained by dividing Rcontrol in equation 3 by the number of bits per symbol, and is in proportion to Rcontrol in equation 3. Accordingly, Rcontrol in equation 3 will be hereinafter referred to as “control information coding rate (coding rate of control information).”\n                                              ⁢                  [          3          ]                                                                              R          control                =                              O                          Q              ′                                =                                    max              (                                                ⌈                                      O                                                                  10                                                                              -                                                          Δ                              offset                              PUSCH                                                                                10                                                                    .                                              R                        data                                                                              ⌉                                ,                                  O                                      4                    ·                                          M                      sc                                                                                  )                        =                          max              ⁡                              (                                                      R                    ⁢                                                                                  ⁢                    1                                    ,                                      R                    ⁢                                                                                  ⁢                    2                                                  )                                                                        (                  Equation          ⁢                                          ⁢          3                )            where ┌x┐ is an integer not greater than x, and max(x,y) is the value of the greater one of x and y. In equation 3, R1 is the coding rate that is set based on user data coding rate Rdata and PUSCH offset, ΔPUSCHoffset, per control information. R2 is the lower limit value of control information coding rate Rcontrol. As is shown in equation 3, control information coding rate Rcontrol is set to the value of the greater value of coding rate R1 and lower limit value R2. A case will be described below where coding rate R1 is greater than lower limit value R2, and control information coding rate Rcontrol is set as coding rate R1.\nIn this case, in equation 3, when PUSCH offset ΔPUSCHoffset is greater than 0, control information coding rate Rcontrol is set lower than user data coding rate Rdata. Generally, unlike user data, control information is not retransmitted. Therefore, by setting PUSCH offset ΔPUSCHoffset, to be greater than 0 and using equation 3, it is possible to lower control information coding rate Rcontrol than user data coding rate Rdata to enhance the capability of error correction of control information."}
{"text": "The present invention relates to benzazabicyclic carbamates of the formulae I and II below, and the pharmaceutically acceptable salts of such compounds. The compounds of formula I are cholinesterase inhibitors and are useful in enhancing memory in patients suffering from Dementia and Alzheimer's disease. The compounds of formula II are novel intermediates used in the synthesis of compounds of formula I. As set forth in detail below, the compounds of formula I and certain of the compounds of formula II are also useful as analgesic agents.\nAlzheimer's disease is associated with degeneration of cholinergic neurons in the basal forebrain that play a fundamental role in cognitive functions, including memory. Becker et al., Drug Development Research, 12, 163-195 (1988). As a result of such degeneration, patients suffering from the disease exhibit a marked reduction in acetylcholine synthesis, choline acetyltransferase activity, acetylcholinesterase activity and choline uptake.\nIt is known that acetylcholinesterase inhibitors are effective in enhancing cholinergic activity and useful in improving the memory of Alzheimer's patients. By inhibiting acetylcholinesterase enzyme, these compounds increase the level of acetylcholine, a neurotransmitter, in the brain and thus enhance memory. Becker et al., supra, report that behavioral changes following cholinesterase inhibition appear to coincide with predicted peak levels of acetylcholine in the brain. They also discuss the efficacy of the three known acetylcholinesterase inhibitors physostigmine, metrifonate, and tetrahydroaminoacridine.\nEuropean Patent 0253372 refers to 1,2,3,3a,8,8-ahexahydro -3a,8 (and 1,3a,8) -di (and tri) methylpyrrolo /2,3-b/indoles of the formula ##STR2## wherein R, R.sub.1, X, Z, and m are as defined in such patent, and states that such compounds inhibit acetylcholinesterase are useful as memory enhancing and analgesic agents.\nEuropean Patent 0154864 refers to physostigmine derivatives of the formula ##STR3## wherein R is (C.sub.2 -C.sub.20) alkyl, branched alkyl, cycloalkyl or aryl, and states that such compounds inhibit acetylcholinesterase and are useful in the treatment of Alzheimer's disease.\nYu et al., Febs Letters, 234, 1, 127-130, (1988), refer to physostigmine derivatives of the formula ##STR4## wherein R is as defined in such article, and discuss their relative potencies as inhibitors of acetylcholinesterase and butylcholinesterase as compared to the corresponding potency of physostigmine.\nAtack et al., J. Pharmacology and Experimental Therapeutics, 249, 1, 194-202 (1989), refer to certain carbamoyl and N(1)-substituted analogs of physostigmine and discuss their relative potency as cholinesterase inhibitors as compared to physostigmine.\nBrufani et al., Eur. J. Blochem, 157, 115-120 (1986), refer to physostigmine analogs of the formula ##STR5## wherein R.sup.1 is an alkyl group and R.sup.2 is hydrogen, and state that such compounds possess anticholinesterase activity.\nKnown cholinesterase inhibitors are useful over a relatively small range of concentrations and exhibit adverse side effects, becoming exremely toxic at concentrations substantially higher than the effective range. Also, the relationship between cholinesterase inhibition and changes in acetylcholine concentrations that follow such inhibition has been shown to be unpredictable, not solely the result of percent cholinesterase inhibition, and strongly affected by the properties of individual drugs. There is therefore a great need for novel cholinesterase inhibitors."}
{"text": "Not applicable.\nNot applicable.\n1. Field of the Invention\nThe present invention relates to the field of fiber optics, and more particularly to the field of optical signal switching, routing, and monitoring.\n2. Description of Related Art\nFIG. 1 is an illustration of an exemplary optical switching system 10 for practicing the invention. For example, optical switching system 10 may represent a 3-dimensional optical switching system. A 3-dimensional optical switching system allows for optical coupling between input fibers and output fibers in different planes using lens arrays and mirror arrays. The lens arrays and mirror arrays provide proper angle and position of light beams traveling from input fibers to output fibers. That is, a light beam must leave and enter a fiber in a direct, beam path.\nReferring to FIG. 1, a generalized optical switching system 10 includes input fiber array 12, first lens array 14, first beam steering apparatus (e.g., Micro-Electro-Mechanical-System, or MEMS, mirror array) 16, second beam steering apparatus (e.g., MEMS mirror array) 18, second lens array 20, and output fiber array 22. System 10 might also include additional or different elements, such as apparatus 24 and 26 for controlling the mirror arrays 16 and 18, respectively.\nInput fiber array 12 provides a plurality of input optical fibers 28 for forming light beams 30 transmitted to (and through) first lens array 14. First lens array 14 includes a plurality of optical lenses 32, which are used to focus beams of light from each input optical fiber 28 to individual mirror devices 34 on mirror array 16. Mirror devices 34 may be electronically, magnetically, or otherwise individually movable to control the beam path of each beam formed by the input optical fibers 28.\nMirror device 34 may be a gimbaled mirror device having a rectangular, elliptical, circular, or other appropriate shape. The plurality of mirror devices 34 for mirror array 16 can pivot a reflective component thereof (not specifically shown in FIG. 1) to redirect or reflect light to varying mirror devices on second mirror array 18. Second mirror array 18 also includes a plurality of mirror devices, similar to those described with regard to first mirror array 16, which are used to redirect and reflect light beams to varying lenses 36 on second lens array 20. second lens array 20 focuses beams of light from second mirror array 18 t o individual output fibers 38 of output fiber array 22.\nOptical switching system 10 allows light beams from any input fiber 28 of input fiber array 12 to be redirected t o any output fiber 38 of output fiber array 22. The above arrangement, including mirror arrays 16, 18 may also be used in scanning systems, printing systems, display systems, and other systems that require redirecting beams 20 of light.\nIt should be noted that for each input optical fiber 28 there is an associated mirror device (such as mirror de vice 34) on mirror array 16, and for each output optical fiber 38 there is an associated mirror device on mirror array 18. In general, there will be a minimum of two input optical fibers and two output optical fibers, and correspondingly two mirrors on each of arrays 16 and 18. There need not be an identical number of input and output optical fibers, although this is typically the case. Also, there will typically be more than two such input and output optical fibers.\nIn general, many types of beam steering arrangements will lend themselves to use with the present invention. For example, two mirror arrays 16, 18 are shown in FIG. 1. Good coupling of a beam into an output optical fiber typically requires controlling fiber position and angle in two dimensions, in addition to the two dimensions of mirror rotation. So, four degrees of freedom are required. However, in some instances, few or greater degrees of freedom may suffice. For example, a one-mirror array shown FIG. 2, where a single mirror array 42 controlled by apparatus 44 is used for directing the beams, or a three or more mirror array arrangement (not shown) may be appropriate. Also, while the above description has been in terms of reflective beam steering, other types of beam steering, such as refractive or diffractive beam steering may employ the present invention with equal advantage (although not otherwise discussed in detail herein).\nCoupling a light beam from an input optical fiber to an output optical fiber requires that the mirrors of mirror arrays 16, 18 be angularly controlled, and that the mirror angles be precise. Such precision, and a degree of control thereover, is needed to optimize output optical power. To obtain such precise positioning, one can either employ or not employ knowledge of actual mirror angle. One method to determine the mirror angles is to add structures to the mirrors that indicate the angle of mirror rotation. One example of this is torsion sensors attached to suspension elements associated with each mirror. See, for example, U.S. Pat. No. 6,044,705. This approach has several disadvantages, including a relatively large number of required interconnections, added processing and manufacturing steps, greater risk of yield losses, etc.\nThe present invention recognizes the limitations of direct measurement of mirror position, and presents an alternative technique for allowing optimization of output optical power. Our approach is to use an optical signal incident on each mirror to determine the mirror\"\"s angle. According to one embodiment, the output optical power is measured, for example at each output optical fiber, and compared with either the input optical power or a reference value, and based on the results of the comparison, a mirror or mirrors are rotated, and output power measured again. This process is referred to herein as feedback stabilization control of mirror position. According to another embodiment, this may be an iterative process, to allow narrowing in on a desired mirror angle for loss optimized output optical power."}
{"text": "This invention relates to interactive wagering systems and methods. More particularly, this invention relates to interactive wagering systems and methods that enable wagers to be placed via a television set-top box, via a computer, and/or via a telephone.\nWagering on sporting events such as horse, dog, and harness racing is a popular leisure activity. However, it is sometimes inconvenient to attend racing events in person. Not all racing fans have sufficient time to visit race tracks as often as they would like and some fans have difficulty in obtaining suitable transportation to the track. Thus, there is a need for wagering services for fans who cannot attend racing events in person.\nOff-track betting establishments, which are generally more readily accessible than race tracks, have attempted to fill this need. However, a racing fan who desires to place a wager still faces the prospect of traveling to the off-track betting establishment.\nWagering via telephone is another option. A user of a telephone-based system typically sets up a telephone account against which wagers may be made. In order to place wagers, the user must interact with a computerized telephone ordering system by pressing appropriate buttons on a touch-tone telephone. This type of system is mainly used for placing wagers. Detailed racing information is typically obtained from other sources, such as printed racing programs.\nWagering via a television set-top box is still another option. As described in Brenner et al. U.S. Pat. No. 5,830,068, which is hereby incorporated by reference herein in its entirety, known systems enable a user to receive wagering information and place wagers using a television set-top box in a user's home.\nWagering via computer is yet another option. Using a computer, a user can similarly receive wagering information and place wagers from the user's home.\nAlthough, systems for telephone, set-top box, and computer wagering are known, no known system provides an integrated wagering system that enables a user to receive wagering information and place wagers using more than one of these methods. Moreover, many known systems for telephone, set-top box, and computer wagering are difficult to use and do not provide a user-friendly interface.\nIt would therefore be desirable to provide systems and methods for interactive wagering that provide an integrated wagering system that can be accessed by telephone, set-top box, and/or computer.\nIt would also be desirable to provide systems and methods for interactive wagering that provide an easy to use interface that allows the user to quickly navigate through the available options to place a wager.\nIt would further be desirable to provide systems and methods for interactive wagering that provide a user interface that provides both an expert wagering mode and a novice wagering mode."}
{"text": "The approaches described in this section could be pursued but are not necessarily approaches that have previously been conceived or pursued. Therefore, unless otherwise indicated, it should not be assumed that any of the approaches described in this section qualify as prior art merely by virtue of their inclusion in this section.\nReliability of Internet access is essential to the ability of households and businesses to manage personal matters and conduct business. A home network of a private customer and a private network of a business enterprise can be connected to the Internet using customer premises equipment, such as routers, modems, switches, and the like. Conventional customer premises equipment collects and stores network-related data, such as settings of the customer premises equipment, setting of electronic devices connected to the Internet using the customer premises equipment, or policies for accessing the Internet by the electronic devices. If the customer premises equipment fails or is replaced with new customer premises equipment, all network-related data stored by the customer premises equipment may be lost and a set-up of the new customer premises equipment may be needed. The customer may need to manually provide settings for the new customer premises equipment.\nIt may be possible for configurations settings associated with customer premises equipment to be stored remotely, for example, on a cloud-based remote storage resource. For example, a virtual servicing node may be in communication with the customer premises equipment and responsible for collecting and storing the network-related data. However, conventional virtual servicing nodes serving the customer premises equipment are not fault tolerant. In case of a failure of the virtual servicing node, the customer premises equipment may lose access to the network-related data and the data connection of the customer network may be interrupted. Thus, failure of a virtual servicing node fail may result in thousands of customers losing data access."}
{"text": "The Molecular Anatomy (MAN) Program at the Holifield National Laboratory is involved in the development of clinical techniques for analyzing biomaterials on a molecular level. The GeMSAEC Fast Analyzer, a type of centrifuge, is one of the results of the program. The centrifuge is capable of analyzing the contents of minute samples placed in separate cuvettes on a rotor. The samples are spun around a light source and the analysis is made using one of four types of measurements, the amount of light the sample (1) absorbs, (2) fluorescences, (3) scatters, or (4) chemiluminescences. The measurements are made by a photomultiplier tube in the analyzer.\nTwo major problems exist with prior art circuits controlling the gain of the photomultiplier. One such prior circuit is described in U.S. Pat. No. 3,783,300, issued Jan. 1, 1974, to Wayne F. Johnson. One problem is the time required for the controller to set the gain on the P.M. tube. It is important to be able to begin the analysis as quickly as possible because reactions can begin to take place in the cuvettes in less than three tenths of a second. The prior method of setting the gain on the P.M. tube takes around 1.5 seconds. By that time, the reaction may have reached a steady state and valuable information is lost. The present invention, to be described hereinbelow, can begin analysis after 3 or 4 revolutions of the rotor. This is approximately 0.2 second at a rotor speed of 1000 rpm.\nThe other problem of prior art systems is that the reference cuvette has to be a water blank due to feedback complications. The prior P.M. tube controller can only use a water blank in the first cuvette as a reference. Other reference solutions tend to have a small degree of instability or noise during a run. A slight variation of the reference cuvette induces a large response in the feedback mechanism of the analyzer. For instance, assume the reference is a fluorescent solution that exhibited a small increase in fluorescence for a short time. The prior control system responds to this small increase by decreasing the results obtained in the cuvettes undergoing analysis. Thus, a small change in the reference solution greatly distorts the analysis results of the other cuvettes.\nThus, there exists a need for an improved system for regulating the gain of the photomultiplier tube of a photometric solution analyzer wherein the gain can be set at a desired or required level and in a shorter time interval, and the analyzer is not limited to a water blank cuvette as a reference. This need has been met in the present invention in a manner to be described hereinbelow."}
{"text": "The use of electronic transmissions has proliferated with the advent of personal computers, or “terminals”, and improved communications networks such as the Internet. Billions of electronic transmissions are sent and received each year in the United States. An electronic transmission, as used herein, is defined as the accessing, processing, or presentation of electronic data, to include word-processed content, mathematical spreadsheets, emails, visual or graphic images, audible content, software code, pattern data, execution commands, computer software programs, Internet web sites, software rule modules, electronic instant messaging, and the like. Such electronic transmissions may take many forms, including: an electronic request for user-customized or user-unique access to stored database content; an electronic request to customize the processing of data according to user-customized or user-unique criteria; and an electronic request to present or display data in a pre-determined, user-customized format.\nIt should be noted that user-customized is different from user-unique. Electronic data or electronic transmissions which are customized to a user, or “user-customized”, have been customized by or for a user, but may not necessarily by unique to that user. Therefore, user-customized data which applies to one user may also apply similarly or identically to another user. However, electronic data or electronic transmissions which are unique to a user, or “user-unique”, are distinctive and without equal, and hence are exclusive to that particular user.\nIn sum, an electronic transmission is the accessing, processing, or presentation of any electronic data or content which does not in and of itself constitute or execute either: an electronic financial transaction wherein the exchange or alteration of any financial assets occurs, nor; an automated door lock or an automated physical site access mechanism.\nA result of the significant popularity of electronic transmissions has been a marketplace transition from using an off-line, individual desktop personal computing model to using an on-line, central-server communications model. Specifically, corporations and individual consumers are moving the main functions of storage, access, processing and presentation of their electronic transmissions from decentralized, unconnected desktop terminals, to centralized databases on servers which service and connect to on-line PCs, known as “client terminals”, via dial-up, leased lines or wireless networks. In this transition, such client terminals are also increasingly being connected to each other. An integrated web of communications is forming that enormously expands the functions and benefits of using such clients, evidenced by fast growth of the Internet and corporate intranets.\nAt the same time, cost reductions in miniaturizing computer hardware components have led to the widespread use of a new generation of computing devices, known herein as “thin-clients”, which are even less expensive and more mobile than traditional desktop terminals and client terminals. The appeal of these new thin-clients is that they offer the potential for the user to send and receive electronic transmissions at virtually any time and from virtually anywhere. Many of these lower cost thin-clients access much of their processing and memory capacities on-line from remote servers via Internet, intranet or extranet connections. These thin-client devices include, but are not limited to: wireless pagers; wireless and tethered telephones; network computers; thin-client exercise machines; electronic books; public access kiosks such as automated teller machines, vending machines, airport information terminals and or public kiosks; hand-held personal digital assistants such as Palm Pilots™ and the like; on-line photocopy machines; automobile embedded Internet-connected appliances which download preferred radio stations, seat and temperature adjustments, and the like; thin-client household appliances such as refrigerators, microwaves, and the like; thin-client home entertainment appliances including on-line televisions such as WebTV™, portable digital audio systems such as the Rio™, along with their associated remote controls.\nThese two trends, of proliferating personal computing devices and of increased on-line communications usage, have led to a distinct problem: with so may personal computing devices, the user now has user-customized electronic data stored on multiple man-made memory devices, or “tokens”, which the user must manage and possess for storage, access, processing and presentation of their electronic transmissions. Further, if the user wants all of these new computer tokens to possess the same capabilities with respect to the user's personalized information and customized functions, then the user needs to frequently and redundantly enter all such user-customized data into each token. This is a cumbersome burden which most consumers eschew. If, on the other hand, the user does not effect such redundancies, then losing or damaging their primary personal computing token would be a severe blow. In this instance, or even in the instance where the user loses or damages a computing token with a subset of their information, then months, and perhaps years, of important personal and likely confidential electronic transmissions could be irretrievably lost, or revealed to an untrusted third-party.\nIn sum, the multitude of such personal computing tokens, whether unconnected desktop terminals or on-line hand held thin clients, has exacerbated the problem of user-reliance on particularly vulnerable, customized memory tokens which can be easily damaged, lost or stolen.\nTo protect these tokens and the resident electronic transmission they contain, the use of various biometrics, such as fingerprints, hand prints, voice prints, retinal images, handwriting samples and the like have been suggested for identification of individuals. However, because the biometrics are generally themselves stored in electronic, and thus reproducible, form on the token itself and because the comparison and verification process is not isolated from the hardware and software directly used by the user attempting access, the problems of fraudulent access and of having to constantly carry these tokens is not alleviated. Further, such systems do not adequately isolate the identity verification process from tampering by someone attempting to gain unauthorized access. Examples of this approach to system security are described in U.S. Pat. No. 4,821,118 to Lafreniere; U.S. Pat. No. 4,993,068 to Piosenka et al.; U.S. Pat. No. 4,995,086 to Lilley et al.; U.S. Pat. No. 5,054,089 to Uchida et al.; U.S. Pat. No. 5,095,194 to Barbanell; U.S. Pat. No. 5,109,427 to Yang; U.S. Pat. No. 5,109,428 to Igaki et al.; U.S. Pat. No. 5,144,680 to Kobayashi et al.; U.S. Pat. No. 5,146,102 to Higuchi et al.; U.S. Pat. No. 5,180,901 to Hiramatsu; U.S. Pat. No. 5,210,588 to Lee; U.S. Pat. No. 5,210,797 to Usui et al.; U.S. Pat. No. 5,222,152 to Fishbine et al.; U.S. Pat. No. 5,230,025 to Fishbine et al.; U.S. Pat. No. 5,241,606 to Horie; U.S. Pat. No. 5,265,162 to Bush et al.; U.S. Pat. No. 5,321,242 to Heath, Jr.; U.S. Pat. No. 5,325,442 to Knapp; U.S. Pat. No. 5,351,303 to Willmore, all of which are incorporated herein by reference.\nAn example of a token-based security system which relies on a biometric of a user can be found in U.S. Pat. No. 5,280,527 to Gullman et al. In Gullman's system, the user must carry and present a credit card sized token (referred to as a biometrics security apparatus) containing a microchip in which is recorded characteristics of the authorized user's voice. In order to initiate the access procedure, the user must insert the token into a terminal such as a public kiosk, and then speak into the terminal to provide a biometrics input for comparison with an authenticated input stored in the microchip of the presented token. The process of identity verification is generally not isolated from potential tampering by one attempting unauthorized access. If a match is found, the remote terminal may then signal the host computer that access should be permitted, or may prompt the user for an additional code, such as a PIN (also stored on the token), before sending the necessary verification signal to the host computer.\nAlthough Gullman's reliance of comparison of stored and input biometrics potentially reduces the risk of unauthorized access as compared to numeric codes, like personal identification numbers, Gullman's use of the token as the repository for the authenticating data combined with Gullman's failure to isolate the identity verification process from the possibility of tampering greatly diminishes any improvement to fraud resistance resulting from the replacement of a numeric code with a biometrics. Further, the system remains cumbersome and inconvenient to use because it too requires the presentation of a personalized memory token in order to initiate an access request.\nAlmost uniformly, prior art disclosing biometrics are token-based systems which teach away from biometrics recognition without user-dependence on personalized memory tokens. Reasons cited for such teachings range from storage requirements for biometrics recognition systems to significant time lapses in identification of a large number of individuals, even for the most powerful computers.\nIn view of the foregoing, there has long been a need for a computerized electronic transmissions system which enables the user to universally access, process and present their electronic transmissions with optimal convenience by not requiring the user to possess any man-made memory tokens on which must be stored the user's customized in order for the user to execute electronic transmissions. Further, there is a need for a tokenless computer system which is highly fraud-resistant, and which is centered around the individual themselves by relying solely upon their unique biometric samples. Such a system should be able to function for the user wherever and whenever the user may be using any generic on-line computing device, whether a desktop or a thin client, for conducting their electronic transmissions.\nFurther, there is a need for a computing system that provides the user with centralized storage, access, processing and presentation of their electronic transmissions regardless of whether the personal computing device the user is using possesses only a resident subset of their user-customized data or in fact possesses none of their user-customized data at all. Further, there is a need for a computerized electronic transmissions system that provides the user with the above benefits whether or not the personal computing device the user may be using at any given time contains powerful resident memory and processing capacities, or whether it contains virtually no resident memory and processing capacities. Further, there is a need for a computer system which relieves the user from having to redundantly data-enter and update a variety of individual personal computing devices in order to achieve the same customized performance from any or all of such devices.\nThere is also a need for a computerized electronic transmissions system which relieves the user from having to redundantly data-enter their personal demographics and customized Internet usage activity information into a variety of Internet web sites in order to achieve uniformly customized service at each such web sites. Additionally, there is a need for a computerized electronic transmissions system which enables a user to benefit from executing customized and complex commands governing their electronic transmissions regardless of whether the on-line computing device the user happens to be using is a high-powered desktop terminal or whether it is a hand-held, ultra thin-client terminal with virtually no resident computer processing or memory capabilities of its own.\nThere is also a need for an electronic transmissions system that uses a strong link to the person being identified, as opposed to merely verifying a user's possession of any physical objects that can be freely transferred.\nThere is a further need for an electronic transmissions system that ensures user convenience by enabling user-authorization without requiring the user to possess, carry, and present one or more proprietary memory tokens, such as man-made user-customized portable memory devices, in order to effect electronic transmissions. Anyone who has lost a smart card or a traditional notebook personal computer, left it at home, had it damaged or stolen knows well the keenly and immediately-felt inconvenience caused by such problems. Therefore, there is a need for an electronic biometric transmissions system that is entirely tokenless.\nThere is another need in the industry for a computerized electronic transmissions system that is sufficiently versatile to accommodate both users who desire to use personal identification codes (PICs), being alphabetical, numerical or graphical, for added security and also consumers who prefer not to use them.\nLastly, such a system must be affordable and flexible enough to be operatively compatible with existing networks having a variety of electronic transmission devices and system configurations."}
{"text": "In LTE-A (Long term evolution advanced) or LTE Rel. 10+, heterogeneous net-work deployments are supported. In a heterogeneous network, a macrocell provides basic coverage, complemented by multiple low-power pico-/femtocells where needed. Such deployments give scenarios of macro-pico/femto inter-cell interferences.\nIn different LTE releases, the inter-cell interference coordination (ICIC) mechanism may be used. In ICIC, the goal is to determine the resources (frequency, time and power) available in each cell, and schedule those resources to users such that the inter-cell interference can be avoided or reduced.\nIn a heterogeneous network (HetNet), base stations (BSs) or evolved Node Bs (eNBs) in terms of 3GPP (3rd Generation Partnership Project) LTE of pico-/femtocells can be operator-deployed or user-deployed. In the operator-deployed case, users of a macrocell are usually allowed to connect to the pico-/femtocell. This is known as open access. And in the user-deployed case, home BSs, femto, or Home eNBs (HeNBs) are installed by the users. Such low-power BSs are typically limited to a small set of user equipment (UE) (e.g. family members) which is called the closed subscriber group (CSG). A UE is usually not allowed to connect to the CSG if it does not have the access code. In this scenario, a non-CSG UE staying in close proximity to the CSG may be subject to strong interference from the HeNB and may not be able to access the macrocell at all.\nIn order to deal with these harsh interference problems, the so-called enhanced inter-cell interference coordination (eICIC) has been introduced, and the further eICIC (FeICIC) may be used as well. By using almost blank subframes (ABSs), the eICIC can schedule the frequency or time resource for macrocell and pico-/femtocell to avoid the interference.\nAdditionally, there is another scenario in which a pico-/femtocell UE is interfered by a strong macrocell. This setup occurs when it is preferable to have a UE connect to a low-power pico-/femtocell node even when it has significantly lower received power as compared to a high-power macro node. This setup may be useful when, for example, the strong cell has weak backhaul quality, or when it is desirable to enable traffic off-loading to the low-power nodes and to achieve true cell-splitting gains in the network. Especially in such HetNet situations, ICIC may not be fully effective in mitigating the inter-cell interference.\nAs ABSs still contain the cell-specific reference signals (CRSs), paging channel (PCH), physical broadcast channel (PBCH), and primary and secondary synchronization channels (PSS/SSS), etc, interference to the transmission at these positions cannot be avoided. Especially due to the large amount of CRSs distributed throughout the overall user bandwidth, which are primarily used for channel estimation, the impact of the interfering CRSs can be quite large. When the CRSs from both cells have different offset, the channel estimation is not affected by the interferer, but the interfering CRSs will wipe out some data tones. In this situation, muting of the interfered data tones is one solution for this problem. When CRSs from both cells have the same or similar offset, they will collide with each other. Based on the colliding CRSs, channel estimation accuracy will degrade severely and the demodulation performance suffers accordingly."}
{"text": "The present invention concerns a device for operating technical medical apparatus, especially dialysis apparatus, with a display screen and a touch screen surface and with means for displaying and/or altering characteristic treatment data. This invention also concerns a method of operating technical medical apparatus.\nDisplay screen units with touch screen surfaces for operation of technical medical apparatus are known. In general, such operating units are used to inquire about actual patient values or machine parameters by touching an appropriate field on the touch screen surface or to enter corresponding setpoints.\nEuropean Patent No. 623,357 describes a device and a method of dialysis where the interface between the apparatus and the user is implemented in the form of such a touch screen monitor. The monitor serves to query and enter characteristic dialysis parameters, such as pumping speed, ultrafiltration rate or the conductivity of the dialysate. To enter a setpoint, the operator must touch a field on the monitor designated with the name of the parameter accordingly, ultimately arriving at a screen view where a corresponding value or a profile can be defined. A generic device and a generic method are disclosed in U.S. Pat. No. 5,609,770, where the dialysis parameters of interest are likewise selected with the help of a touch screen monitor. After the selection is made, a desired setpoint can be assigned to the parameters, with the selection of parameters as well as the inputting of values taking place by touching the surface accordingly. Because of the multitude of variable dialysis parameters, a display screen view contains a correspondingly large number of selectable fields, which makes rapid and reliable operation of the apparatus difficult. Working without mistakes is also impaired by the fact that the fields are all the same size and can be differentiated only on the basis of their designation. It is thus possible to touch the wrong fields on the touch screen surface, especially when working under time pressure, and thus, for example, enter a parameter setpoint incorrectly, which can have considerable consequences in terms of risk for the patient during treatment.\nThe object of the present invention is to simplify the handling of technical medical apparatus and reduce the probability of operating errors.\nStarting from a generic device and a generic method, the object of the invention is achieved by a device for operating a technical medical apparatus comprising a display screen and a touch screen surface and a first means for displaying and/or altering characteristic treatment data. In addition, a second means are provided for display of symbols characteristic of the components of the apparatus, where the first means can be controlled by touching the symbols on the touch screen surface. This makes it possible to largely rule out any confusion of symbols due to inadequate differentiation of the fields of the touch screen surface to be touched. Instead, by simply touching the respective symbol of the component to be influenced, the display of the corresponding data, e.g., in a control field, is activated. Such a graphical user surface permits safe and reliable use of a dialysis machine, for example. The medical personnel will very quickly become familiar with the scope of the device according to this invention due to these characteristic symbols, thus reducing the familiarization period and also increasing reliability in operation.\nIt is especially advantageous if at least one complete dialysis cycle can be displayed with the symbols for the components contained in the cycle. Either the dialysis fluid cycle or the blood cycle or both may be displayed, with the components contained therein, such as pumps, cut-off devices, sensors or the dialysis machine itself can be represented by means of characteristic symbols. Displaying the entire cycle facilitates the user\"\"s orientation in selecting one or more components of the dialysis cycle. This minimizes the probability of confusing different pumps, cut-off devices or heating devices, for example. It is also possible that only a portion of the components of a medical technical apparatus, in particular a dialysis apparatus or a dialysis circulatory system, can be depicted. In this context, two or even more device components of any kind, as well as their functional interrelationship, can be depicted. For example, it is conceivable that the dialyzer and the blood pump, as well as the hose line connecting the two components can be indicated using the second means. The selection of the combinations of components or segments to be indicated from the respective entire installation can be made according to the basic settings of the apparatus, or they can be set by the operator according to need. By depicting the functional connections of the interrelationships, the mode of operation for these components is made easier.\nIt is also possible that only portion of the components of a medical technical apparatus, in particular a dialysis apparatus or a dialysis circulatory system, can be depicted. In this context, two or even more device components of any kind, as well as their functional interrelationship, can be depicted. For example, it is conceivable that the dialyzer and the blood pump, as well as the hose line connecting the two components can be indicated using the second means. The selection of the combinations of components or segments to be indicated from the respective entire installation can be made according to the basic settings of the apparatus, or they can be set by the operator according to need. By depicting the functional connections of the interrelationships, the mode of operation for these components is made easier.\nIn another embodiment of this invention, the first and second means are designed so that the symbols can be displayed simultaneously and the respective treatment data can be displayed and/or altered. This makes it possible to avoid accidentally changing a setpoint or entering an actual value because the symbol of the component itself and the parameter to be changed are always displayed simultaneously.\nThe treatment data that can be displayed and/or changed by the first means may include the designation as well as setpoints and/or actual values and/or limit values of the instrument parameters and/or patient parameters. Whereas in a monitoring mode, essentially the actual values of the characteristic data of the components are queried, but in another operating mode of the device according to this invention, specific changes can be made in setpoints. This makes it possible to define not only individual setpoints but also a time dependence of the setpoints and thus a setpoint profile. Likewise, it is possible to set limit values for the characteristic data which, if maintained, will reliably prevent a hazardous operating state. If the limit values are reached, this may be signaled visually or acoustically or it may lead to a shutdown of the apparatus or cause cut-off devices to close. The data that can be displayed and/or changed by the first means may also include other parameters and values concerning, for example, the course of the dialysis treatment so far or the storable program data from preceding treatments of various patients.\nIt is especially advantageous if the means for changing the treatment data and the means for displaying the treatment data can be controlled simultaneously by touching the symbols. This generates control surfaces, for example, which display the current actual values of the corresponding parameters simultaneously and also include, for example, fields marked with arrows for corresponding changes in setpoints. Changing the data includes not only changing the setpoint but also complete startup or shutdown of a component or a function unit of the apparatus, for example.\nIt is likewise possible for the treatment data to be displayed first after touching the characteristic symbols for the components of the medical apparatus and for the means for altering the data to be controllable by appropriate touch of the displayable treatment data. Thus, after touching a symbol once, first only the data are displayed, and it is necessary to touch the touch screen monitor again only in the event the data must be changed.\nThe components of the technical medical apparatus that can be displayed by means of symbols may include blood and dialysate pumps and/or sensors and/or means for processing the dialysate. The sensors may be designed as both temperature and pressure pickups. It is likewise possible for the dialysis machine itself or cut-off devices such as valves or clamps or even heating devices to be displayable.\nIt is especially advantageous if several of the respective treatment data can be displayed and/or changed simultaneously by touching multiple symbols. In this way it is possible, for example, to obtain a set of actual values or a list of setpoints by touching several symbols or all symbols. It is likewise possible in this way to obtain a profile of the actual value trends in the treatment administered previously or other past treatments.\nIn another embodiment of this invention, the first means are designed so that a time dependence of the treatment data can be defined. In this way it is possible to define setpoint profiles, so that the success of the treatment can usually be improved in comparison with fixed predetermined setpoints within the context of a dialysis treatment. The setpoints, for example, can be entered accordingly with the time stipulations in the form of a table or a similar list. It is also possible for the user to predefine a desired course of the setpoint as a function of time by touching a corresponding field on the touch screen surface. In this case a computer unit determines the desired setpoints to be reached after a preselectable dialysis time.\nIt is especially advantageous if the touch screen function can be turned off after a period of time has elapsed after manipulating or touching the touch screen surface. Consequently, the touch screen function is switched off when a defined period of time has elapsed since the last operation or the last manipulation of the touch screen surface, thereby preventing any negative effects due to inadvertent touching. In particular, there can then no longer be any inadvertent change in display or even a change in the setpoints for patient parameters or instrument parameters. It is also possible for a display comprising a selection of characteristic data to be generated simultaneously or after a time lag. Thus, a survey display of only relevant data can be generated, making it easy for the operator to monitor a treatment process and clearly detect malfunctions. The selection of characteristic data that can be displayed in this way can be varied by the user or may be predetermined by the apparatus and depend, for example, on the type of treatment. An especially simple display is achieved when characteristic data are selected in an enlarged display.\nThe present invention also concerns a method of operating a technical medical apparatus, in particular a dialysis machine. The characteristic treatment data can be displayed and/or altered by a display screen with a touch screen surface in a first screen view, and in at least one second display screen view, the symbols characteristic of the components of the dialysis machine are displayed. The first display screen view is generated by touching these symbols. The method according to the present invention thus requires only that the user touch the desired symbol of the component of which either the actual value or another parameter is to be called up or a setpoint is to be changed. After touching the characteristic symbol, the desired actual values or other data are displayed in a corresponding screen view. Instead of displaying the data, it is likewise possible for a change in the data to be enabled. This is the goal in particular when the actual values of the parameters are always displayed together with the symbols of the components of the apparatus and when renewed display of these values is not necessary. After displaying or altering the parameters, the first screen view may be closed, and the symbols of the components of the apparatus are displayed again, for example.\nIn another embodiment of this invention, a screen display is generated by touching several symbols simultaneously or in succession, thus permitting simultaneous display or alteration of multiple treatment data. Thus, for example, the actual values of all treatment parameters may be displayed simultaneously and thus can be monitored easily.\nAccording to another embodiment of the present invention, touching the symbols generates the display of the corresponding treatment data, and touching the displayed treatment data generates another screen display which makes it possible to change these data.\nIt is especially advantageous if the symbols of the components of the apparatus are always displayed simultaneously with the respective actual values of the parameters. Such an embodiment of the method permits an overview of the current data without requiring intervention by the treating personnel. In this case it is sufficient if touching the corresponding symbols does not lead to renewed display of the actual value but instead permits a change in the corresponding setpoint or other data, for example.\nIn another embodiment of the present invention, touching the symbols generates a screen view by means of which the time dependence of one or more of the treatment data is displayed. For example, if it is necessary for the ion concentration of a dialysis solution to be variable over time, the time dependence of this setpoint can also be entered by touching a corresponding symbol."}
{"text": "Technical Field\nEmbodiments of the invention relate generally to time synchronization of networked components. Particular embodiments relate to time synchronization of CT scanner subsystems and computer equipment.\nDiscussion of Art\nCT scanners are used in medicine both for diagnosis and for guidance of interventions. In a CT scanner, an X-ray source and a detector (such as an array of scintillation crystals with associated photomultipliers or photodiodes) are positioned opposite each other on a gantry, which rotates around a target such as a patient. The X-ray source is activated to emit X-rays, some of which pass through the target and others are reflected or absorbed by the target. X-rays that pass through the target to the detector generate signals.\nThe signals at the detector vary temporally due to stochastic fluctuations in absorption and reflection of the X-rays, and also vary spatially as the source and detector rotate around the stationary target, which thereby presents different absorption and reflection characteristics at each angular position of the gantry. The variations of the detector signals, during one or more rotations of the gantry, can be mathematically correlated or “computed” in a central processing unit, which generates correlation data suitable for producing a cross-sectional image representation of different types of tissue or other material within the target.\nComputer processors are known for implementing software, e.g., the algorithms necessary for operation of a computed tomography (CT) scanner components to emit and detect X-rays and to correlate detections with gantry positions and with source position relative to the target. Computer processors implement software instructions by switching electronic logic gates to manipulate binary HIGH or LOW signals (“bits”). Computer processors can be implemented in RISC, ASIC, FPGA, or other architectures.\nAccuracy of correlation, and, consequently, image quality, depend to a great extent upon accurate logging of the X-ray source and detector positions relative to the target. Typically, these positions are calculated based on elapsed time from a zero-time reference position, and based on measurements of gantry angular velocity derived from, e.g., a time integral of gantry motor current, voltage, or shaft speed. Alternatively, positions may be directly measured from a bar code or the like marked on a surface of the gantry. In any case, each detector signal and each position reading typically is tagged with the local time at the detector or position sensor, and for enhanced quality, the X-ray source power level also may be continuously or periodically measured and tagged with the local time at the source. Thus, timekeeping is a critical aspect of CT scanner design.\nIn complex systems like CT scanners, multiple processors are networked together. For example, a CT scanner may include an X-ray source processor, an X-ray detector processor, a position sensor processor, and a gantry control board processor, all of these being configured for communication with each other via a network. Typically, networked processors communicate with each other by sending data packets. For coordinated action, networked processors attempt to maintain synchrony so that each received data packet is implemented or passed on in proper sequence. Each processor has its own clock (comprising a time counter driven by an oscillator, e.g., a bistable), and synchrony can be maintained by periodically transmitting and receiving a time signal, from a designated “master” clock to the “slave” clocks of the other processors. The component clocks periodically adjust their time counters in response to the time signals received from the master clock. Between time signals, the component clocks increment their time counters based on the pulses produced by their respective oscillators.\nThe master clock/time signal paradigm is often employed because it is data-efficient and has low network overhead. Typical periodicities for master clock time signals range from about 1 MHz to about 200 MHz and each processor switches its logic gates anywhere from about 1000 times to about 5 times between receipt of each time signal from the master clock. As a result, a large quantity of process data can be generated and transmitted between time signals, using the same signal path used for the time signals (“path sharing”). By contrast, continuous synchronous network time protocols (continuous time signal transmissions) are not favored because such protocols require a dedicated physical layer that cannot also be used for transmitting process data.\nA potential issue with the time signal paradigm is that, during the numerous switching operations between time signals while the component clocks increment their respective time counters, asynchrony or clock “drift” can emerge between the master clock and each component clock and also among the component clocks. The clock drift phenomenon arises because each component clock oscillator has its own characteristic frequency that infinitesimally varies from the frequencies of the other clock oscillators. Over an extended sequence of switching operations, the different frequencies of the different oscillators can result in significant divergence of the component clock time counters.\nAnother potential issue with the time signal paradigm is that a unique network transit time or signal delay exists between the master clock and each component clock. Accordingly, each component clock receives the time signal at a different time relative to the master clock. This time difference can change according to environmental factors that can vary the transmission delay. Some possibly relevant environmental factors include, for example, thermal strain, changes in impedance, electromagnetic field interactions. The ordinary expected result of this variation in transit time is that each component clock time “jitters” around the master clock time.\nFor example, image quality of a CT scanner can be limited by “clock drift” or “jitter” between the processors associated with the detector, the source, the position sensor, and the gantry control board. In particular, variance of detector local time away from source or position sensor local time, or from gantry control board local time, can result in erroneous position readings, which limit the achievable quality of the correlation data. Depending on the error threshold designed into the particular CT scanner, small clock drifts can degrade image quality without establishing an error condition.\nIEEE 1588 improves on the time signal paradigm by implementing a Best Master Clock (BMC) algorithm that is run by all of the processors, in order to select which processor will run the master clock. One feature of the BMC algorithm is that the processors cooperatively select a master clock in a manner that is not necessarily determinative from a user standpoint—in other words, the BMC algorithm is not intended to “force” selection of any particular processor as the master clock.\nIEEE 1588 addresses an issue of component clock drift due to component clock oscillator frequency differing from master clock oscillator frequency by instituting a reciprocal measure of time signal delay. The reciprocal measure of time signal delay is used by each component clock for adjusting its local time from the time of receiving the master clock pulse, to more closely approximate the simultaneous local time at the master clock. Determining the reciprocal measure of time signal delay, for each component clock, requires repeated reciprocal communication at scheduled intervals between the master clock and the component clocks."}
{"text": "1. Field of the Invention\nThe invention relates to a feed-forward structure, as compared to a feed-back structure, for canceling the dc-offset of a down-converted RF signal in a Direct Conversion Receiver (DCR).\n2. Description of the Related Art\nFor the purpose of indicating the background of the invention and/or the state of the art, the following four references are incorporated by reference in their entirety: B. Razavi, xe2x80x9cDesign Considerations for Direct-Conversion Receivers,xe2x80x9d IEEE Transactions on Circuits and Systems, vol 44, No. 6, June 1997 (herein xe2x80x9cReference 1xe2x80x9d); W. Namgoong, xe2x80x9cDirect-Conversion RF Receiver Designxe2x80x9d, IEEE Transactions on Communications, vol. 0.49, No. 3, March 2001 (herein xe2x80x9cReference 2xe2x80x9d); A. Abidi, xe2x80x9cDirect-Conversion Radio Transceivers for Digital Communicationsxe2x80x9d, IEEE Journal of Solid-State Circuits, vol. 30, No. 12, December 1995 (herein xe2x80x9cReference 3xe2x80x9d); and U.S. Pat. No. 6,317,064 to Ferrer, et al., xe2x80x9cDC offset correction adaptable to multiple requirements,xe2x80x9d published Nov. 13, 2001 (herein xe2x80x9cReference 4xe2x80x9d).\nWireless systems have become an essential part of modern life. Systems such as cell phones, wireless LANs, etc., are very popular and the penetration rate of these technologies is increasing very fast throughout the world. In view of high demand and marketplace competition, emphasis is placed upon new systems having small size, low power consumption, low cost, and of course, high bandwidth and quality service.\nOne of the impediments to making wireless equipment less expensive, smaller, and lower power, is the difficulty of integrating the analog front-end. The heterodyne receiver is the most popular architecture for wireless systems\"\" front-end. It has good selectivity and performance but suffers from bulky off-chip components that can not be integrated with current technologies.\nAn alternative solution is the direct conversion receiver (DCR), which has a very simple architecture and a potential for a fully integrated front-end, with low cost, small size, and low power consumption. Such features are particularly attractive for next generation wireless handsets and for software defined radio (SDR) systems.\nHowever, DCR suffers from some serious problems like dc-offset, 1/f noise, I/Q mismatch, and even-order distortion. I/Q mismatch and even-order distortion can be made negligible with good circuit design techniques, but dc-offset and 1/f noise are generally more serious and challenging problems, dc offset being considered the most challenging problem in realization of DCR in commercial products. Strong dc-offset drives the baseband amplifiers to nonlinear mode, and for ideal amplifiers it saturates analog-to-digital converters (ADCs) (e.g., Reference 1).\nA very low frequency and reasonably high offset voltage may appear at the output of the mixer in a DCR. This offset voltage dominates the signal strength by a factor of 50-100 times in amplitude and highly degrades the bit-error-rate probability if not removed. Unless unwanted offset is removed in the analog domain prior to sampling and before the baseband amplifiers, the offset saturates the baseband amplifiers and causes devastating non-linear distortion. Also in the case of ideal amplifiers, a very large dynamic range ADC is needed to resolve the signal from the dc-offset (e.g., References 1 and 2).\nFIG. 1A demonstrates a conventional DCR in which dc-offset cancellation circuits 20 remove, in the analog domain, dc-offset from the down-converted signal prior to variable gain amplifiers (VGA) 7.\nThere are different sources that cause dc-offset in a DCR. The local oscillator (LO) 4 signal leaks to the radio frequency (RF) port and the mixer 3 down converts the leaked signal to baseband, i.e., a zero intermediate frequency (IF). The LO signal also leaks from antenna 1 and reflects off external objects and self-down converts to direct current. These self-mixings introduce tremendous dc-offset because the LO signal is generally much stronger than RF signal. In addition, a strong interferer may leak to the LO port of the mixer and be translated to zero IF. The amount of dc-offset is difficult to predict because its magnitude changes with receiver location, orientation, and time. Therefore in addition to dc-offset estimation, a tracking method should be incorporated in the receiver to follow the offset variations (e.g., References 2 and 3).\nA simple and straightforward method to cancel dc-offset is ac-coupling, as demonstrated in FIG. 1B. This method is suitable and cost effective where the signal spectrum does not have so much energy at dc and the ac-coupling does not degrade system performance. A pager is a good example for this case. Additionally, for Code Division Muliple Access (CDMA) systems, especially Wideband CDMA, some degradation in the spectrum is affordable and ac-coupling may be a viable solution. However, for CDMA, the corner frequency must be in the range of few kHz, requiring large off-chip capacitors, slowing response time, which is not desirable. None of these drawbacks are favorable in making a fully integrated front-end for sophisticated wireless systems.\nAnother approach to canceling dc-offset utilizes digital, rather than analog, techniques. Most digital methods for dc-offset cancellation are feedback-based average-and-subtract techniques. They estimate the offset noise by averaging the digitized signal after the entire analog baseband section. Then the estimated dc-offset value fed back and subtracted from the down-converted signal at the mixer output via adder 11.\nAn example of this technique is illustrated in FIG. 2. The signal is sampled after the ADC 9, but before the digital baseband receiver 10. The sampled signal is digitally processed by a dc-offset cancellation circuit 30 containing an averaging circuit 31, memory 32, and a digital-to-analog converter (DAC) 33. The dc-offset cancellation circuit 30 estimates the offset noise, which is then subtracted from the down-converted signal at adder 11.\nDifferent methods have been proposed for estimating dc-offset, but nearly all of them share some drawbacks, such as (1) slow response due to their feedback structure, (2) saturation of the ADC and baseband amplifiers due to passing offset noise through the entire analog baseband section during detection time, (3) need for high resolution ADCs, (4) high complexity and computationally inefficiency, and (5) inefficient dc-offset tracking methods.\nTo overcome the drawbacks of the prior art, a new feed-forward dc-offset structure and method is proposed. The new canceller comprises a feed-forward structure where dc-offset, including offset noise, is estimated from down-converted signal, and cancelled before applying the signal to the analog baseband section. DC-offset is cancelled without using feedback.\nA direct conversion receiver, according to an exemplary embodiment of the present invention, comprises a dc-offset tracking unit, a subtraction circuit, and an analog baseband circuit. The dc-offset tracking unit samples a down-converted signal, and outputs a dc-offset signal based upon the dc-offset of the sampled signal. The dc-offset tracking unit comprises a linear digital filter, estimating and tracking dc-offset in the sampled down-converted signal. The subtraction circuit receives the down-converted signal after the down-converted signal is sampled by the dc-offset tracking unit, and subtracts the dc-offset signal from the received down-converted signal. The analog baseband circuit comprises a variable gain amplifier, and receives the down-converted signal after the dc-offset has been subtracted.\nPreferably, the linear digital filter is a first order recursive filter, having a transfer function:                               y          ⁡                      [            n            ]                          =                              x            ⁡                          [              n              ]                                +                      ky            ⁡                          [                              n                -                1                            ]                                                                                    y            ⁡                          [              n              ]                                -                      ky            ⁡                          [                              n                -                1                            ]                                      =                  x          ⁡                      [            n            ]                                                            H          ⁡                      (            z            )                          =                                            Y              ⁡                              (                z                )                                                    X              ⁡                              (                z                )                                              =                      1                          1              -                              kz                                  -                  1                                                                        \nk being a loop gain of the first order recursive filter.\nThe direct conversion receiver further includes a mixer, an analog-to-digital converter, a low pass filter, and a digital baseband receiver. The mixer receives an RF signal and down-converts it to zero IF, outputting a down-converted signal. The analog-to-digital converter receives the down-converted signal after the dc-offset is subtracted by the subtraction circuit and after the gain of the signal is adjusted by the analog baseband circuit, converting the down-converted signal from analog into digital, for processing by the digital baseband receiver. The low pass filter filters the down-converted signal after the down-converted signal is output by the mixer, but before the signal is converted to digital by said analog-to-digital converter.\nThe dc-offset tracking unit further comprises an ADC and a DAC. The ADC samples the down-converted signal, and the DAC outputs the dc-offset signal, the linear digital filter being connected in the digital domain between the ADC and DAC. Instead of a DAC, a digitally programmable dc-voltage source can be used.\nThe dc-offset tracking unit further comprises a sampling rate converter, connected between said linear digital filter and the DAC.\nThe DCR further comprises an anti-aliasing filter, filtering the sampled signal prior to the signal being input into the ADC of the dc-offset tracking unit.\nThe invention further comprises a method of canceling dc-offset from a signal. An exemplary embodiment comprises monitoring a signal, detecting the dc-offset in the monitored signal by applying a linear filtering function, and subtracting the detected dc-offset from the signal at a point after where the signal was monitored.\nThe method preferably comprises using a first order recursive function as the linear filtering function, the first order recursive function having a transfer function of:                               y          ⁡                      [            n            ]                          =                              x            ⁡                          [              n              ]                                +                      ky            ⁡                          [                              n                -                1                            ]                                                                                    y            ⁡                          [              n              ]                                -                      ky            ⁡                          [                              n                -                1                            ]                                      =                  x          ⁡                      [            n            ]                                                            H          ⁡                      (            z            )                          =                                            Y              ⁡                              (                z                )                                                    X              ⁡                              (                z                )                                              =                      1                          1              -                              kz                                  -                  1                                                                        \nk being a loop gain of the first order recursive function.\nThe invention further comprises a method operating a direct conversion receiver. An exemplary embodiment comprises receiving an RF signal, down-converting the RF signal to zero IF by mixing the RF signal with a signal from a local oscillator, canceling a dc-offset in the down-converted signal, adjusting gain of the down-converted signal after the dc-offset is cancelled, converting the gain-adjusted signal from analog to digital, low pass filtering the down-converted signal, and processing the digitally-converted signal. Canceling the dc-offset in the down-converted signal comprises monitoring the down-converted signal, detecting the dc-offset in the monitored signal by applying a linear filtering function, and subtracting the detected dc-offset from the down-converted signal at a point in the direct conversion receiver after where the down-converted signal was monitored. Low pass filtering is performed after the RF signal is down-converted, but before the gain-adjusted signal is converted from analog to digital.\nPreferably, the linear filtering function is a first order recursive function, having a transfer function of:                               y          ⁡                      [            n            ]                          =                              x            ⁡                          [              n              ]                                +                      ky            ⁡                          [                              n                -                1                            ]                                                                                    y            ⁡                          [              n              ]                                -                      ky            ⁡                          [                              n                -                1                            ]                                      =                  x          ⁡                      [            n            ]                                                            H          ⁡                      (            z            )                          =                                            Y              ⁡                              (                z                )                                                    X              ⁡                              (                z                )                                              =                      1                          1              -                              kz                                  -                  1                                                                        \nk being a loop gain of the first order recursive function.\nDetecting the dc-offset further comprises converting the monitored signal from analog to digital before applying the linear filtering function, and converting the dc-offset detected by the linear filtering function from digital to analog.\nThe method further comprises anti-alias filtering the monitored signal prior to detecting the dc-offset.\nConverting the dc-offset from digital to analog further comprises adjusting a digital sampling rate.\nThese structures and methods have very low complexity and are computationally efficient and accurate, and are very suitable for making low cost, small size, and low power consumption integrated RF receivers based on DCR.\nThe feed-forward design has many advantages over conventional designs, such as fast response time, more relaxed requirements on the analog baseband section, no need for high-resolution A/D converters to digitize a desired signal, negligible low frequency spectrum distortion that occurs in some of the wireless standards, very good tracking capability, easy configuration for different operating situations, suitability for multimode/multi-standard systems, no need for memory cells to store sample signals, and computational efficiency and less complexity because of the ability to operate at lower sampling rates and lower resolution."}
{"text": "A microcontroller unit (MCU) is small computer formed on an integrated circuit. MCUs can generate signals for controlling of a wide range of devices, such as electric motors, voltage regulators, office machines, appliances, implantable medical devices, etc.\nMCUs typically include a central processing unit (CPU), memory, and programmable peripherals components. The CPU executes a program (hereinafter referred to as an embedded program), which is typically stored in flash memory. MCUs may also include analog-to-digital converters (ADCs), digital-to-analog converters (DACs), comparators, timer/counter channels (hereinafter timer channels), etc. Timer channels are often used to autonomously control external devices such as induction motors. Because the timer channels operate autonomously, the CPU can perform other functions while timer channels control external devices."}
{"text": "The rapid acceptance of forming hay into large cylindrical bales has not been matched with complementary development of means for the lifting, loading, or transportation of the bale so produced. This is of particular consequence because, depending upon the type of grass baled, cylindrical bales may weight between 1000 and 2000 pounds green and remain even as heavy as 1500 pounds when dry. Further, medium sized cylindrical bales typically range approximately five feet in diameter and approximately five feet in length. At this bulk and weight, prior practices of loading the baled hay by hand onto trucks or trailers are no longer practical.\nThe most popular expedient making use of common, currently available equipment is to employ a front end loader attachment mounted on a tractor. However, such systems require their own independent hydraulic and control systems, are time consuming to both mount and remove, and create frequent maintenance demands. Further drawbacks to the use of front end loaders as applied to loading and unloading of hay bales are that they are not standard equipment on the majority of farms and that a special procurement towards this use proves an expensive solution, both as to purchase and maintenance. Furthermore, agricultural tractors are not ordinarily designed to be heavily loaded at the front; they are usually much stronger at the rear.\nIn addition, several special handling devices have been developed, but none combines the simplicity, convenience and versatility necessary to effect a satisfactory solution.\nThe Nelson hay loading device disclosed in U.S. Pat. No. 2,954,886 is a fairly complex tractor mounted apparatus which teaches the use of gripping jaws to engage a bale and a hydraulically actuated swinging extension to rotate the bale from a position beside the tractor to an elevated position from which the bale may be dropped onto an elevated surface behind the tractor. The apparatus is not suitable for conveniently unloading hay bales from an elevated surface, nor for applications requiring more careful placement than the \"dropping\" operation allows. The operation contemplated by Nelson requires an independent hydraulic system operating from a rigid, stationary hitch and the resulting design is not adaptable to a configuration that is readily attachable to a standard, hydraulically actuated three point hitch.\nCox, in U.S. Pat. Nos. 4,099,629 and 4,015,739, discloses a hay bale handling machine, including one configuration to be mounted on a tractor. His primary disclosure is of a device that may be mounted into the bed of a pickup truck. This pickup truck apparatus uses an extended spear to impale a bale of hay and then pivots the immediate upright member from which the spear protrudes until the spear is essentially vertical at which point the bale is in a cylindrically upright position, resting over the truck bed upon the previously upright and now horizontal spear carring member. This configuration does not allow either loading or unloading hay bales upon any appreciably elevated surface. The disclosures also contain a configuration for mounting the spear upon a three point hitch behind a tractor, but again the device is not appropriate to work with elevated storage or transportation surfaces. The spear is mounted upon an upright post which is directly attached to the three point hitch system. The top of the post is attached through a top support arm to the upper hitch point of the tractor. The bottom of the post is shown permanently affixed to a drawbar disposed between the lower hitch members. Both the top support arm and the drawbar are standard means for attaching implements to a three point hitch and neither is directly adaptable for mounting the post and spear in a configuration that allows loading or unloading from elevated surfaces. Rather it is the object of Cox to engage and tilt the bale for transportation only upon the spear or spike itself, and not for loading to or unloading from elevated surfaces.\nThe identical failings are present in the Jones disclosure for U.S. Pat. No. 4,120,405. Indeed, here the inadequacies of the three point hitch configuration seems to be conceded as this system is discussed anciliary to a front end loader adaptation for handling hay bales. But even as adapted, the front end loader retains the inadequacies discussed in the preceding for standard, commercially available front end loaders.\nThe tractor mounted embodiment of McFarland's U.S. Pat. No. 4,084,707 is primarily of interest in unrolling hay bales, though it may be used for transport in the same sense as the Cox device. However, it too is fundamentally lacking for the immediate application.\nFurther, while various elevating extension schemes have been developed to be mounted upon a tractor's three point hitch, none is appropriate for the contemplated use.\nA shovel device is disclosed by U.S. Pat. No. 4,068,774, however Howell's apparatus is not adaptable because it teaches a single extension bar device. Such a one-bar system could not produce an appropriate tilt of an engagement spear.\nThe Larsen device, under U.S. Pat. No. 2,311,671 shows a manure loader. This device includes complications which were appropriate to facilitate a dumping action, but are not appropriate for a hay bale loading machine where dumping is not contemplated. Further, the structural bar system itself disposes a lower bar so low as to unduly limit an approach to an elevated platform.\nThe Holopainen ditch digging device disclosed in U.S. Pat. No. 2,496,874 is an inappropriate elevating extension in the present application for the same reasons as Larsen. In addition, the Holopainen system calls for auxiliary hydraulic systems."}
{"text": "The present invention relates to a high efficiency cardiac valve structure, which substantially consists of an annular portion having, at one axial end thereof, a pair of valve members of substantial arc-shaped configurations, which valve members are hinged or pivoted to the annular body.\nMore specifically, said valve members are so arranged and designed as to provide a one-way valve able of opening as the blood flows in a direction, and of automatically closing as the blood flows in the opposite direction.\nCardiac valves are generally known which, even if they have been found to be satisfactory from the operational standpoint, are however susceptible to further improvements, mainly with respect to the coupling of the valve members to the ring element."}
{"text": "Separation techniques using open tubes, such as Capillary Chromatography and Capillary Electrophoresis (CE) have become universal and important tools for separating inorganic anions, peptides and proteins. This is because of the high separation efficiency, the good resolution, the faster separation, and the lower consumption of samples and solutions that can generally be achieved compared to liquid chromatographic techniques.\nAlthough CE has been known for several decades, new developments and applications continue to be forthcoming with approximately 3,000 papers presently published annually in the field. (Hu et al, Anal. Chem. 2002, 74, 2833-2850.) One area of research within the field of CE is the study of wall chemistry. This research focuses on ways to prevent wall adsorption and suppress EOF, which in turn improve the separation efficiency and resolution of different analytes. For example, inorganic anions are very difficult to separate with native silica capillaries unless very high salt concentrations and very low pH are used to suppress the EOF, because the EOF direction is opposite to the electrophoretic movement of anions. For anions with high electrophoretic mobility, the separation time becomes much longer, whereas anions having low electrophoretic mobility may in extreme cases be transported out of the capillary by the EOF at the injection end of the capillary. The negative charge that the native silica capillary surface expresses over a wide pH range is also problematic in the separation of positive charge molecules such as basic proteins because of electrostatic interaction. Therefore, the EOF must be suppressed or even reversed by shielding or altering the surface charge, in order to obtain fast separations and improve separation efficiency. Many approaches, including adjusting the concentration or pH of the buffer, adding small cationic-type solutes or various detergents or polymers, and attaching a static covalently-bonded polymer layer on the lumen surface have been developed to fulfill this task. (Rodrigueza et al, Anal. Chem. Acta 1999, 383, 1-26; Righetti et al, Electrophoresis 2001, 22, 603-611; Liu, Electrophoresis 2001, 22, 612-628; Horvath et al, Electrophoresis 2001, 22, 644-655; Chen et al, Anal. Chem. 1993, 65, 2770-2775; Pietryk et al, J. Chromatogr., A 1997, 775, 327-338; Minnoor et al, J. Chromatogr., A 2000, 884, 297-309.) All these modifications have shown varying degrees of improvement on the separation of different samples.\nAlthough zwitterionic separation materials, prepared either by dynamically coating or covalently bonding zwitterionic functionalities, have been used for long time in liquid chromatography (Hu et al, P.R. Trends Anal. Chem. 1998, 17, 73-79; Nesterenko et al, Anal. Sci. 2000, 16, 565-574; Jiang et al, Anal. Chem. 1999, 71, 333-344; Anal. Chem. 2001, 73, 1993-2003; Anal. Chem. 2002, 74, 4682-4687; Viklund et al, Macromolecules 2000, 33, 2539-2544; Viklund et al, Anal. Chem. 2001, 73, 444-452.), there are only limited studies involving modified capillaries in CE with zwitterionic functionality. Most of these studies are focused on suppressing the EOF and preventing wall adsorption by zwitterionic surfactants. The different approaches have demonstrated a potential for improving the separation efficiency of inorganic and organic anions or cations, peptides, and proteins. (U.S. Pat. No. 5,415,747; Swedberg, J. Chromatogr. 1990, 503, 449-452; Strege et al, J. Liq. Chromatogr. 1993, 16, 51-68; Greve et al, J. Chromatogr., A 1994, 680, 15-24; Gong et al, Electrophoresis 1997, 18, 732-735; Yeung et al, Anal. Chem. 1997, 69, 3435-3441; Baryla et al, Anal. Chem. 2000, 72, 2280-2284; Woodland et al, Analyst 2001, 126, 28-32; Castelletti et al, J. Chromatogr., A 2000, 894, 281-289; Yokoyama et al, Anal. Chem. 2001, 371, 502-506; Mori et al, Anal. Bioanal. Chem., 2002, 372, 181-186.) Nonetheless, the dynamic coating procedure requires that the zwitterionic surfactant be added to the separation solution in order to get a reproducible surface, a prerequisite for good reproducibility in the separation.\nThe addition of zwitterionic surfactant results a more complicated separation matrix; their presence causes interference in certain detection schemes such as mass spectrometry (Baryla et al, J. Chromatogr., A 2002, 956, 271-277; Cunliffe et al, Anal. Chem. 2002, 74, 776-783.) and also adds to the complexity of samples that are fraction collected for subsequent use or analysis. Although Cunliffe et al. have tried to use the zwitterionic surfactant with double carbon chains to obtain semi-permanent layers on the capillary surface, the stability of the layer is still an important issue. This is because the mixture containing the zwitterionic surfactant and other chemicals has to be used to flush the capillary between runs of separating proteins, which results longer experiment time. (Cunliffe et al, Anal. Chem. 2002, 74, 776-783.) Accordingly, there is a need in the art for synthesis of capillaries with covalently-bonded zwitterionic functionalities both in order to investigate basic separation properties and for comparison with capillaries dynamically coated with zwitterionic surfactants. U.S. Pat. No. 5,415,747 suggests that a zwitterionic species may be bonded to capillaries though a direct silicon-carbon linkage. However, no results are provided regarding a capillary which comprises covalently bonded zwitterionic groups. Consequently, there is a need for further improvements in the field of capillary electrophoresis and capillary chromatography."}
{"text": "Prostaglandins are a family of biologically active lipid acids that possess as a common feature the prostan-1-oic acid structure ##STR2## Prostaglandins are grouped into the types E, F, A, B, C and D, based on the presence or absence of certain functionalities in the cyclopentane ring. The numerical subscripts, as in, for example, prostaglandin \"E.sub.1 \" and prostaglandin \"E.sub.2, \" refer to the number of unsaturated bonds in the side chains; the subscripts \".alpha.\" or \".beta.\" as in prostaglandin \"F.sub.1.alpha. \" or prostaglandin \"F.sub.1.beta.,\" refer to the configuration of substituents in the ring.\nBiological activities of the prostaglandins include stimulation of smooth muscle, dilation of small arteries, bronchial dilation, lowering of blood pressure, inhibition of gastric secretion, lipolysis, and platelet aggregation, induction of labor, abortion, and menstruation, and increase in ocular pressure. PGE.sub.1, specifically, is known as a bronchodilator and a vasodilator, and is also known to stimulate release of erythropoietin from the renal cortex and to inhibit allergic responses and blood-platelet aggregation. PGE.sub.2, the most common and most biologically potent of the mammalian prostaglandins, acts to contract the uterine muscle, to inhibit gastric acid secretion and protect the gastric mucosal lining and, like PGE.sub.1, has been established to be a bronchodilator and a stimulant of the release of erythropoietin from the renal cortex. Properties of the various prostaglandins have been reviewed extensively; see, e.g., Ramwell et al., Nature 221:1251 (1969).\nBiosynthesis of the prostaglandins occurs by enzymatic conversion of unsaturated twenty-carbon fatty acids. For example, the endoperoxides PGG.sub.2 and PGH.sub.2 are prepared by the action of the enzyme complex prostaglandin cyclooxygenase on the lipid precursor arachidonic acid, while PGE.sub.1 is biosynthesized by enzymatic conversion of 8,11,14-eicosatrienoic acid. Biosynthetic studies of the prostaglandins are reviewed in Samuelsson, Prog. Biochem. Pharmacol. 5:109 (1969).\nA number of synthetic routes to the various prostaglandins have been explored. A twenty-step synthesis to PGE.sub.2 and PGF.sub.2.alpha., starting with thallous cyclopentadiene, was developed by Corey et al. The process has also proved to be useful in the synthesis of the one-series prostaglandins. Catalytic reduction of bis-protected PGF2.alpha. results in selective saturation of the 5,6-double bond to afford ##STR3## followed by transformations leading to PGE.sub.1 and PGF.sub.1.alpha.. See, e.g., Corey et al., J. Am. Chem. Soc. 91:5675 (1969), Corey et al., J. Am. Chem. Soc. 92:2586 (1970) , and Corey et al., Tetrahedron Letters 307 (1970).\nAnother synthetic route uses norbornadiene as a starting material, while still another approach involves the use of racemic bicyclo[3.2.0]hept-2-en-6-one as a starting material. This latter synthesis involves an enantioconvergent approach, such that the following two enantiomers are used, obviating the need for optical resolution: ##STR4##\nHowever, there remains a need in the art for a simpler and more direct route to the various prostaglandins, particularly prostaglandin E.sub.1, prostaglandin E.sub.2, and derivatives thereof. It is also desirable that such a synthesis provide the desired product in relatively high yield, be simple and straightforward to scale up, and involve inexpensive, commercially available reagents. The present invention is directed to such a method, and involves the use of a lithium copper reagent in a \"one-pot\" synthesis of PGE.sub.1, PGE.sub.2 and analogs thereof."}
{"text": "1. Field of the Invention\nThis invention relates to a device and method for treating ingrown nails on toes and fingers.\n2. Background of the Related Art\nA normal nail on a toe or finger has a relatively large radius or curvature corresponding generally to the circumference of the toe or finger. However, certain conditions may cause the edge or edges of the nail to become curved or curled inwardly, especially at or near the end of the nail. As a result, the interned edge or edges of the nail protrude into the flesh of the toe or finger, giving rise to the condition known as ingrown toenail or ingrown finger nail.\nSince an ingrown nail may cause discomfort and in some cases severe pain, a number of devices and appliances have been devised to treat the condition. For example, U.S. Pat. No. 137,106 to E. E. Stedman discloses a curved spring which overlies the toe and engages the edges of the nail. U.S. Pat. No. 1,596,532 to F. Haener discloses a double walled guard with portions adapted to lie above and below the nail, respectively, and double walled arms extending rearwardly to embrace the edges of the nail. U.S. Pat. No. 2,342,530 to E. F. Coates discloses a protective nail device which includes two complementary sections, with each section having a thickened edge portion which lies along the edge of the nail and a relatively thin upper portion which lies along the top of the nail, and the device extends forward of the end of the nail to hold the flesh in its proper position with respect to the nail. U.S. Pat. No. 2,505,086 to R. P. Andrews discloses a wide sheet of plastic with an attaching flap which is inserted under the ingrown side of the nail. U.S. Pat. No. 2,920,621 discloses a nylon channel inserted under the edge of the nail to serve as an abutment for a lipped metal strap which lies horizontally across the face of the nail and is then taped into position with adhesive tape. U.S. Pat. No. 3,981,298 to Philip G. Vironda discloses a liquid material spread across the base of the nail, causing the material to harden and to become adhered to the nail and to serve as a support to retain the form of the nail as it grows. U.S. Pat. No. 4,674,486 to Ronald G. Hoffman discloses a resilient strip of material placed on the upper surface of the toenail, which urges the side edges of the nail upwardly and thereby tending to fatten the toenail and eliminated the inwardly curved edges of the ingrown nail. U.S. Pat. No. 5,012,799 to Werner G. Remmen discloses a steel wire spring having a central inverted V-shaped bridge, and the free ends of the wire are curved downwardly to form small lips slipped over the side edges of the toenail.\nA number of nail trimming devices and instruments have also been devised for the purpose of cutting and treating ingrown toenails, including U.S. Pat. No. 2,479,514 to L. C. Rucker, U.S. Pat. No. 3,799,160 to Harry L. Hahn, U.S. Pat. No. 4,936,322 to Damian G. DeSantis, and U.S. Pat. No. 4,964,213 to Patricia A. Suggs.\nDespite the number of devices and methods which have been developed for the treatment of ingrown nails, none have been fully satisfactory to the podiatry and beauty industries. A significant number of people continue to have discomfort and severe pain as a result of ingrown nails."}
{"text": "1. Field of the Invention\nThis invention relates to memory circuits, and more particularly to a memory cells formed from floating-body transistors and their applications.\n2. Description of the Related Art\nIntegrating memory circuits (e.g., caches) on microprocessor dies is a power efficient way of achieving higher performance. In fabricating these dies, a direct correlation exists between memory cell density and performance; that is, performance generally increases with the number of cells in the memory. The degree of integration that can be achieved is not limitless and most often is constrained by memory type, the limited die space available for accommodating the memory circuits, and manufacturing costs as well as other practical considerations.\nOne memory type that has been proposed for integration onto a microprocessor die is SRAM memory. When implemented in large caches, each SRAM memory cell uses six transistors to read and write the data therein. This six-transistor configuration consumes considerable die space and thus leads to low memory cell density. This, in turn, places a significant limitation on microprocessor system performance, and for at least these reasons using SRAM cells to fabricate on-die caches is considered to be less than optimal.\nAnother memory type proposed for integration onto a microprocessor die is DRAM memory. As shown in FIG. 1, a conventional DRAM cell uses one transistor T and one capacitor C to control the storage of data therein. Switching of the transistor is controlled by the voltage state of a word line WL and data is conveyed along a bit line BL via a sense amplifier. Because fewer components are used, an on-die cache formed from this so-called 1T/1C structure has far greater cell density and data storage capacity than SRAM memory. In fact, DRAM cells can be ten times smaller in area compared with SRAM cells. For this reason, DRAM has been favored by chip designers over other types of memory for forming caches on microprocessor dies.\nWhile they are favored, conventional DRAM caches are not without drawbacks. For example, the switching transistor and storage capacitor in each cell are implemented as separate elements. This necessarily reduces the cell density of the memory, which has an adverse affect on system integration and performance. Also, costly processing steps are required to make a capacitor in each cell that can store enough charge to maintain reasonable refresh times (e.g., typically at least 25 ff is needed).\nA need therefore exists for a memory cell which has a higher density than other cells which have been proposed and which therefore may be relied on to improve, for example, the on-die integration of caches in microprocessor systems. A need also exists to construct a memory cell of this type at low cost using existing manufacturing technologies."}
{"text": "The following description is provided to assist the understanding of the reader. None of the information provided or references cited is admitted to be prior art.\nIn the normal course of business, a company can utilize a number of vendors to procure needed products and services. As a company can be both a buyer and a vendor, a buyer refers to a company that is purchasing goods or services, and a vendor is a company that sells goods or services. FIG. 1 is a graph 100 of a common relationship between a buyer's amount of purchases and the number of vendors used by the buyer. Some buyers spend a large amount at a small number of vendors (102). These buyers though tend to still have numerous vendors that provide smaller amounts of products and services (104).\nIn instances where the buyer is a very large corporation, with very substantial financial resources, and the vendors that represent the majority of purchases are also very large corporations, with very substantial financial resources, costly, complex and customized electronic supply chain order, fulfillment, invoice presentment and payment systems are sometimes implemented. Once these complex systems have been implemented, the large buyer and vendor are able to interact with one another in a nearly paperless fashion, including purchase order delivery, invoicing and payment. These systems are often referred to as closed loop environments, as the platforms are not extendable, without deep and costly integration, to all other customers or vendors. Often, these systems rely on electronic data interchange (EDI) interfaces, provided by third parties, to deliver data to one another. And, when payment is due, these systems rely on integrations with bank software platforms to deliver electronic payment via the Automated Clearing House (ACH). The deep customized integration that can exist in these highly complex and costly systems, by its very nature, requires substantial effort, often completed by third party consultants, to design and implement. And, once the system is designed, both buyer and vendor must still allocate substantial management and staff resources during the setup process. All of the above ignores completely that finding vendors and buyers, completing and responding to requests for proposals (RFPs), and “setting up vendors and customers” within each other's accounting or enterprise resource planning systems is time consuming, manual and relies heavily on the acquisition of information that can only be provided directly by the counter party. That said, once the system is complete, efficiencies still improve. In other circumstances such deep integration is simply not practical. For instance, in the retail sector, consumer tastes change so rapidly that many purchases from vendors must be managed on a moment-by-moment basis. Notwithstanding, large buyers with substantial financial resources still work very hard to eke out as much efficiency from the procurement process as possible, often at the expense of their much smaller vendors, who are forced to comply with the variant invoicing requirements and standards imposed upon them by their large customers. And even then these systems are generally designed for vendors that supply valuable inventory, and are not implemented with small infrequently used vendors supplying items that would generally be classified as general and administrative expense items.\nIn many cases, these large buyers eschew deep integration and simply require their vendors to communicate with them concerning transactional information via Electronic Data Interchange. In these cases, the vendor is generally required to hire a third party integrator to serve as an intermediary to enable them to receive each of their customers' electronic orders and comply with each of their customers' unique invoicing requirements. The introduction of a third party integrator is an expensive necessity for vendors who choose to do business with these customers. It should be noted that it is highly unusual that all of a vendor's customers require EDI integration, and as a result, vendors are required to maintain many disparate processes to support and manage their many customer relationships. Some customers receive invoices via EDI, some via paper, others via email. Managing the funds owed by their customers is equally inefficient for small and large vendors alike.\nIn rare circumstances, large buyers with substantial financial resources have gone to the trouble and expense of developing, or purchasing, Electronic Invoice Presentment and Payment (“EIPP”) solutions to communicate invoice and payment status to vendors. However, this remains rare, and even where it exists, the invoice and payment information is buyer specific. That is to say, when vendors sell to buyers that have a platform that provides invoice and payment information, each system is unique and independent, a closed loop. And even then, there is no standardization within these platforms (the information available is different, both in terms of its type, presentation, access, etc.), and as a result, even in the rare circumstances where an online system is available, the acquisition of invoice and payment status across customers is time consuming and inefficient. In most cases, no data is available online and vendors must call each of their customers on the phone, or email them, to obtain invoice and payment information, necessitating customer response to these calls and emails.\nAll of the above consume substantial time and effort for both the buyer and vendor.\nElectronic invoice and presentment solutions are generally designed to meaningfully improve the efficiency of a large buyer with a large number of disparate vendors. This improvement is accomplished in numerous ways, but heavily focuses on the elimination of paper handling, data entry, and manual invoice approval among accounts payable departments. The solutions are, by design, buyer-centric. The buyer receives increased efficiency at the expense of the vendor, who must comply with specific invoice requirements mandated by the buyer. This arrangement often means that the vendor must submit invoices to the buyer following very precise business rules that often force vendors to ignore the invoice printing and email functions that exist within their accounting software platforms, instead submitting invoice information to various customers using unique buyer-specific one off processes. While electronic invoice presentment and payment systems can integrate into a buyer's accounting system, they rarely integrate with similar efficiency into the systems of their vendors, and then, only within custom designed and highly complex systems available only to the largest vendors with substantial financial resources. Even then, the integration is unique to the specific buyer-vendor relationship, rather than to all buyers, and conversely, all vendors. In the rare cases where a customer has made an online invoice status system uniquely available to the vendor, if a vendor needed to request the modification of an invoice after delivery to their buyer, the vendor would need to make the change within their own accounting system, and then resubmit the corrected invoice, or credit memo, to their customer, in a duplicative process.\nSmaller vendors with fewer financial resources, or that sell small quantities of goods or services, particularly on an infrequent basis, or larger vendors with substantial financial resources that sell a small amount of product to their customers, simply have too little capital, sophistication, or financial incentive to undertake such integration projects or implement an EDI relationship. In these circumstances, relationships generally involve the production and mailing of a large number of paper invoices, manual buyer approvals, data entry, collection phone calls and, ultimately, payments by check."}
{"text": "In the case of adjustable steering columns for motor vehicles, the steering spindle carrying the steering wheel and thus also the steering wheel can be adjusted relative to a mounting part by means of which the steering column can be secured on the body of the motor vehicle, this representing a comfort function of the steering column. Conventional adjustment options here are length adjustment and/or height or rake adjustment. In addition to adjustable steering columns, in which adjustment after opening a clamping mechanism is carried out manually, electrically adjustable steering columns are also known. With these, adjustment in the adjustment direction for length adjustment and/or in the adjustment direction for height or rake adjustment is accomplished in a motorized manner by means of at least one drive motor. In this case, a conventional embodiment envisages that the drive motor should be part of a spindle drive. Spindle drives, which are also referred to as worm gear screw mechanisms or screw link actuators, have a threaded spindle on the external thread of which a threaded nut is arranged in order to convert a rotary motion into a translational motion. Either the threaded nut is driven via appropriate mechanism elements by the drive motor, and the threaded spindle is adjusted axially by the rotation of the threaded nut, or the threaded spindle is driven and the spindle nut is adjusted axially by the rotation of the threaded spindle. The threaded nut or the threaded spindle is generally driven by means of a worm arranged on a motor shaft, which is in engagement with worm toothing arranged on the outer circumference of the spindle nut (in the case where the spindle nut is driven) or on the outer circumference of a gearwheel rigidly connected to the threaded spindle (in the case where the threaded spindle is driven).\nAn electrically adjustable steering column having a spindle drive is disclosed in DE 10 2007 039 361 A1, for example.\nOne general problem with electrically adjustable steering columns is that of reducing noise during adjustment. For this purpose, DE 10 2006 036 183 B4 proposes a belt drive, which separates the electric drive motor acoustically in terms of drive from the adjusting mechanism.\nDE 10 2009 057 388 proposes a piezoelectrically driven active vibration damping system. Both technologies are very complex and expensive.\nIn the case of the electrically adjustable steering column known from EP 1 486 395 B1, a transmission part with a spherical head projecting into a cylindrical sleeve is provided on the axially adjustable spindle nut of the spindle drive. This sleeve is provided with an elastomeric lining, whereby noise reduction is also achieved. This technology is also complex since the guide by means of which the spindle nut transmits the motion to the adjustable steering column part has to be of very precise design in order to avoid rattling and jamming. The noise absorption achieved by means of the embodiment shown in this publication is also limited.\nWO 2010/069434 A1 discloses embodiments of electrically adjustable steering columns, wherein, for the purpose of noise reduction, a housing part, in particular that of an electric drive motor, and/or a holding part is/are produced from a sheet embodied as composite material, wherein the composite material has a first metal layer and a second metal layer, between which a layer composed of a viscoelastic material or of an elastomer is arranged."}
{"text": "The invention relates to a composite plate, especially for raised floors to be supported at the corners, consisting of a shell which is open at the top, manufactured of tension-resistant material, preferably sheet steel, and a filling of pressure-resistant material, of which the density and strength can vary within a relatively wide range, and also with anchoring means in the shell for generating a binding effect between the filling and the shell and a reinforcement connected with the bottom plate of the shell.\nA composite plate which is similar to the above structural plate is known from German Patent 20 04 101 and has a shell of sheet steel forming the outside wrapper of the plate and also has a pressure-resistant material embedded in the shell, and the pressure-resistant material is generally concrete or anhydrite. This shell is produced by the deep-drawing manufacturing process, the result of which is that the material of the side walls which are drawn upward during manufacture is less thick than the material of the bottom plate of the shell. Because of this more than anything else, the bearing strength or respectively the carrying capacity of such a composite plate is considerably greater in the middle than at its border. The danger therefore arises that when this composite plate is carrying a load, for instance, when it is travelling with a heavy truck filled with files, the border of the plate becomes gradually further deflected than the middle of the plate. Thus a step occurs in the configuration of the plate and with each impact of the truck of files against said step there occurs one more instance of dynamic stress, which with repetition can lead to permanent deformation of the composite plate or even to destruction of the material (formation of tears in the shell and/or the filling).\nThe drawbacks described above also arise with composite plates in which the thickness of the material of the side walls is identical to that of the floor plate of the shell and/or in which, as in German Offenlegungsschrift 24 45 854, reinforcement rods are welded to the floor plate of the shell adjacent to the side walls. With the just aforementioned composite plates, too, upon application of progressively higher stress the deflection at the border of the plate is greater than the middle of the plate, since the moments of inertia and resistance on the border of the plate are not sufficiently heightened as a result of the presence of reinforcement rods welded to the floor plate of the shell.\nWith known composite plate embodiments, in which the armoring traditionally is embedded within the concrete or the like (cf. for instance German Patent 26 16 317), in comparison with the composite plate of the presently described structural type with a shell-shaped outside wrapper, the light construction characteristics, in other words the ratio between the carrying capacity of the composite plate of Haus and its own weight, is considerably less favorable, so that simply because of these factors these composite plates do not even come into consideration in this case."}
{"text": "(a) Technical Field\nThe present disclosure relates to an apparatus and a method for controlling a vehicle equipped with a motor, and more particularly, to a control apparatus and a control method for enhancing drivability in an electric vehicle capable of recovering electric energy from kinetic energy of the vehicle through the motor during coasting of the vehicle and capable of charging a battery.\n(b) Background Art\nGenerally, an electric vehicle refers to an automobile that is driven by means of electric power, i.e., an automobile that is driven by means of power of a motor and equipped with a motor as a source of power.\nAn electric vehicle is an environmentally-friendly vehicle with less pollutant emissions than traditional vehicles equipped with an internal combustion engine (ICE) using fossil fuel such as gasoline or diesel fuel.\nExamples of the electric vehicle include a pure electric vehicle (EV) that is driven by a motor powered by electric power stored in a battery, a hybrid electric vehicle (HEV) that is driven by both a motor and a combustion engine, and a fuel cell electric vehicle (FCEV) that is driven by a motor powered by electric power generated from a fuel cell.\nAmong them, the pure electric vehicle and the fuel cell electric vehicle are driven only by motor power, while the hybrid vehicle is driven by motor power or combined power of the motor and the internal combustion engine (hereinafter, this is simply referred to as “engine”).\nThe disclosure of this section is to provide background of the invention. Applicant notes that this section may contain information available before this application. However, by providing this section, Applicant does not admit that any information contained in this section constitutes prior art."}
{"text": "This invention relates in general to ram jet rocket engines and in particular to a new and useful control valve for regulating the flow of precombustion gases between a precombustion chamber and a main combustion chamber.\nIn machine construction, control valves for metering the flow of fluids and hydraulically or pneumatically operated devices are known in a variety of designs. In general, such movable mechanical component parts are required to operate exactly, comply with the operating conditions, and be adequately shaped to minimize losses in flow. Further, these control members are to be as lightweight as possible, to reduce mass forces and permit high operating frequencies.\nBall controlled valves are also used in various machines. For example, German OS 28 16 806 shows in FIG. 5 a ball valve in an antiblocking device for vehicle brakes, in which the ball controlling the flow aperture is actuated by a push rod. The fluid flows through the valve, both in front of and behind the control aperture, in the axial direction.\nWith certain modes of operation, particular difficulties arise which call for specific properties of the employed valves, to overcome the problems. For example, in some thermal machines, such as ram jet rocket engines, operating with fuel-rich gases which are produced in a precombustion chamber from a solid chemical fuel and flow through one or more conduits into a main combustion chamber to which air oxygen is supplied for final stoichiometric combustion, solid particles are obtained in the hot gas stream. Being adhesive, these particles are extremely troublesome, they soil lines and, particularly, form deposits on the edges, in corners, recesses and cavities of the valves which may cause failures of these important elements.\nThe prior art valve construction of the above mentioned German OS 28 16 806 is still not capable of safely eliminating these damaging phenomena, since no particular measures are taken to actively prevent or fight deposits in the critical spaces, i.e. in the ball cavity and behind."}
{"text": "The present invention pertains to the field of buses. More particularly, the present invention pertains to clock generation and recovery for a bus.\nMany systems have a bus including a set of transmission lines on which information and one or more clock signals are communicated between various circuits. For example, digital computers have a bus for communicating among circuits such as processors, memory, direct memory access controllers, graphics processors among many other circuits. The bus may include transmission lines for data, addresses, and a system clock. The clock signal on the system clock line provides a time reference by which circuits connected to the bus are synchronized. Keeping an accurate time reference influences the reliable transfer of information on the data and address lines and the reliable use of the bus by the various circuits connected to the bus.\nUnfortunately, conventional clock signals may be 1) susceptible to many sources of error and 2) may contribute to the noise environment of the system. For example, most clock signals have a single-ended substantially square or trapezoidal shape that is characterized by a fundamental frequency (Fo), duty cycle, and edge rate. The nature and characteristics of the clock signals make them susceptible to a variety of undesirable effects produced by impedance discontinuities and noise, among other sources of error. Furthermore, the broadband or wide band characteristics of conventional clock signals may cause interference with other signals and devices in the system.\nFirst, because of the shape of the signal, the signal typically has significant high frequency components which are higher in frequency than the fundamental frequency. Consequently, the signal may occupy a bandwidth that is five times to ten times the fundamental frequency. Unfortunately, transmission lines do not typically provide uniform transmission of the wide band of frequencies contained in the typical clock signal. For example, any type of impedance discontinuity is frequency dependent, causing signal integrity problems that add skew to the system. This problem may result in the improper latching of data causing improper operation.\nSecond, the clock signal is affected by skew/jitter that may result from power plane noise, simultaneous switching output noise coupling, and impedance mismatches in the transmission lines. The effects of these skew/jitter sources are difficult to reject at a receiver that receives a single-ended signal. The effects of these skew/jitter sources on the signal may lead to improper operation and lack of synchronization.\nThird, conventional clock signals, such as substantially trapezoidal or rectangular waves, present a significant electromagnetic interference (EMI) problem because of the many harmonics contained in such signals. The presence of many harmonics with significant power creates fields that may undesirably couple various elements in a system or interfere with other signals in the system.\nSince having a wide band clock signal causes problems in crossing impedance discontinuities, it would be advantageous to generate a clock signal that is less problematic when crossing impedance discontinuities. Additionally, since generating a single-ended signal is not effective for minimizing the effects of the skew/jitter sources, it would be advantageous to generate and recover a clock signal in a manner that would allow rejection of the effects of the skew/jitter sources. Furthermore, since the presence of many harmonics in a clock signal is undesirable, it would be advantageous to generate a clock signal which does not present a significant EMI problem.\nAccording to an embodiment of the invention, a method for recovering a clock signal communicated on a system bus is described. The method includes receiving at a receiver a first signal having a first polarity and receiving at the receiver a second signal having an opposite polarity to the first signal. The method also includes generating at the receiver a first clock signal based upon the first signal and the second signal."}
{"text": "Gas turbine engine controllers are designed to satisfy both performance and safety requirements. To achieve this, engine control logic includes three primary functions: steady-state control, transient control, and limit protection. The steady-state control logic is designed to maintain output power at a demanded level in the presence of disturbances and uncertainties. The transient control logic is designed to safely transition the engine from one power level to another within performance requirements. The limit protection logic is designed to ensure that critical parameters never violate constraints. These functions are all integrated via loop selection logic.\nIn a typical power management loop, the power request is converted to a setpoint. The difference between the setpoint and the corresponding feedback signal is used to drive power management compensation, which ultimately produces a rate command that is output to the loop selection logic. The loop selection logic either passes the rate command through to a common integrator or selects a rate command that is appropriate for the current engine operating mode.\nIn order to develop control logic for steady-state control, transient control, and limit protection, linear point models may be used to approximate the gas turbine engine system. Specifically, control laws are developed and tuned in a computational environment in which linear point models are used to represent the engine dynamics. Development and tuning occurs for one linear point model at a time, and the resulting individual control laws are integrated together to form control logic satisfying both performance and operability requirements across the complete operating envelope of the engine. The controller is then tested on a non-linear engine system model where additional tuning is accomplished. Once modeling and simulation development is completed, bench and engine testing is conducted to fine tune and verify the control logic. All of the models used for control design generally represent a single engine condition.\nIf the control logic is to perform acceptably, it must safely accommodate nonlinearities and uncertainties not present in the design models. In addition, control logic must accommodate the effects of disturbances, manufacturing variations, and degradations that will be present in all of the components of the system. Conventional control logic accommodates these factors by implementing margins within the control logic design. However, this approach may reduce engine performance by not allowing the control logic to access engine states in these margins."}
{"text": "The present invention relates to a microcomputer system with a watchdog timer, and in particular to such a microcomputer system with watchdog timer which incorporates a safety enhancing system for detecting faults in the hardware of the system.\nIn a conventional microcomputer system, a watchdog timer of a per se known sort is provided externally to the microcomputer, and the clock signal source for the watchdog timer circuit is provided separately from the system clock of the microcomputer, so that, in the event of any failure in the clock signal source for the system clock, the watchdog timer circuit may be able to detect any failure in the hardware of the microcomputer. However, in such a per se conventional system, because the clock signal source of the watchdog timer circuit is independent from the system clock signal source, in the event of a failure of the clock signal source and/or the counter of the watchdog timer circuit, it is impossible to detect the fact from the outside of the system, i.e., by the user. As a result, even if the clock signal source and the counter unit of the watchdog timer circuit should fail, the microcomputer keeps on operating independently from it, and therefore, even if the microcomputer should run wild, i.e. the system should run into a certain locked cycle or should not behave normally, due to any failure in the system, the watchdog timer circuit cannot stop it.\nTherefore, in Japanese Patent Application No. 59-78061, which was filed on Apr. 18, 1984, and which it is not hereby intended to admit as prior art to this application except to the extent otherwise required by law, the assignee of the present patent application, as Japanese applicant, proposed to include a process in the system program of the microcomputer for monitoring the timer clock of the watchdog timer circuit and in case of any abnormal behavior thereof for producing an output to the outside indicating such an occurrence.\nWith such a construction, in the case of any failure in the watchdog timer circuit, the microcomputer can detect it, and the safety in case of any failure in the hardware of the system may be enhanced.\nHowever, according to such a construction, a monitoring clock source is provided for the watchdog timer separately from the clock source from the system clock, and there is a problem that, should this monitoring clock source fail, even when the system action of the microcomputer is normal, it becomes impossible to detect any abnormality of the watchdog timer circuit."}
{"text": "The present invention relates to the safe disposal of needles which may be contaminated by blood or other fluids. The needles for disposal are attached to syringes using rotatable connectors, typically, but not limited to, threaded luer lock or luer slip type. The present invention is filed the same day, has the same inventors as and is assigned to the same party as SYSTEM FOR FACILITATING THE REMOVAL AND SAFE DISPOSITION OF MEDICAL NEEDLES. That patent application, Ser. No. 08/235,388 with a filing date of Apr. 29, 1994, is the invention of Nelson, Vangeison, Evans and Wells and is incorporated by reference in this application as if it were fully set forth herein.\nThe SYSTEM FOR FACILITATING THE REMOVAL AND SAFE DISPOSITION OF MEDICAL NEEDLES invention relates to a system and structure for facilitating the removal of potentially contaminated needles and depositing them into a container in an operation which requires the operator to use only one hand and virtually no physical effort. The system and structure employ at least one electrically powered and controlled rotating device for removing the needles. The system uses feedback control sub-systems for operating and controlling the rotating device, for limiting the number of needles which can be placed in a single storage container and for positioning the needle receiver in a predetermined configuration. The present invention teaches a novel and inventive safe for mating with the system and forming a part of that system as well as being inventive in its own right.\nThe basic problem to be solved in the field is the safe disposal of used (possibly bloodborne pathogen contaminated) needles with as close to zero risk of needle contact with personnel who use the needles as is possible. Greater risk to personnel is present when the needle removal device requires two hands or one hand with significant physical force or significant manual dexerity such as when reinserting the needle into its protective sheath. Greater risk to personnel is also present when the needles are not placed in a tamper proof container or when the container can overfill. A device which signals its full condition is safer to users than one which is merely in a jammed (inoperative) condition without otherwise signalling its jammed or full status."}
{"text": "Arabinoxylan, a polysaccharide composed of xylose and arabinose, is part of the water soluble and insoluble fibre present in cereals, in particular in the cell walls. Hydrolysis of arabinoxylan is an important prerequisite for improved utilization of cereal hemicellulose, e.g. in the ethanol fermentation industry and other cereal-based industries.\nArabinoxylan consist of alpha-L-arabinofuranose residues attached as branch-points to a beta-(1-4)-linked xylose polymeric backbone. The xylose residues may be mono-substituted in the C2- or C3-position or di-substituted at both the C2- and C3-position. In addition, ferulic acid and p-coumaric acid may be covalently linked to arabinoxylan via esterification at the C5 position of some of the arabinosyl units. These substitutions on the xylan backbone retard the actions of xylanases and the complete hydrolysis of arabinoxylan thus requires both side-group cleaving and depolymerising activities. The major products of hydrolysis of arabinoxylan are the C5 sugars xylose and arabinose.\nA process for hydrolysis of arabinoxylan using synergistic interactions amongst enzymes present commercial enzyme compositions from Humicola insolens and Trichoderma reesei has previously been described by the present inventors in Sørensen, H. R. et al. (Biotechnology and Bioengineering, Vol. 81, No. 6, March 20, 726-731, 2003). Enzyme catalyzed hydrolysis of >50% of the soluble part of the wheat endosperm arabinoxylan could be achieved, but only low monosaccharide yields were obtained with similar enzymatic treatments on insoluble wheat arabinoxylan. However, since the arabinoxylan degrading enzyme activities are present as side-activities in commercial preparations having other enzyme activities as their main activity, high dosage levels of 5-10 wt % of the enzyme preparation per weight of the substrate has to be added for obtaining efficient hydrolysis. Such high enzyme addition levels are not feasible for use in full scale production applications and improved processes for hydrolysis of arabinoxylan are thus needed."}
{"text": "A pedestrian crosswalk, which is a traffic marking painted on a road surface, is a place designated for pedestrians to cross a road. A pedestrian crosswalk is often found at intersections or the middle sections of roads where many pedestrians cross. Today, as most of the pedestrian crosswalks, a zebra crossing having alternating dark and light stripes is used for reminding car drivers to give way to the pedestrians crossing the road. However, the use of such traffic markings is not enough for controlling and directing traffic in a place where the traffic volume of pedestrians and vehicles is high, and there are often accidents caused by pedestrians and vehicles struggling for smooth traffic on the same road. Therefore, traffic signals, also known as traffic lights, are also provided at the pedestrian crosswalk to control flows of traffic. However, such traffic signals only provide light signs for pedestrians to follow when crossing the road, but are unkindly for a visually impaired person.\nIn recent years, due to the society's concern for the safety of vulnerable groups, access-free facilities have also become a part of urban construction. Therefore, in order to provide an accessible environment for visually impaired persons in pedestrian crosswalks, audible traffic signals are also provided. The audible traffic signals, when changed to green, indicate a walking direction of pedestrians by producing different sounds to help visually impaired persons recognize the light signs for easily crossing the pedestrian crosswalk. However, although such audible traffic signals can produce different sounds to indicate different directions to guide visually impaired persons, the sounds are emitted in all directions and are easy to cause confusion, which makes the visually impaired persons unable to recognize the walking direction, and deviate from the walking direction when crossing, resulting in danger or accident. Furthermore, the visually impaired persons have to be trained to be able to recognize the directions represented by different sounds of the audible traffic signal, and it is very inconvenient for them. Moreover, in order for the visually impaired persons to clearly hear the sounds of the audible traffic signal, the sound volume must be greater than the ambient noise. However, such a sound volume is nothing but a noise that affects the lives of nearby residents and also causes them physical and mental harm after a long-term exposure."}
{"text": "Data centres that provide an easily accessible store of data are a necessary resource for many entities. The amount of data stored in a data centre may be in the order of multiple petabytes. It is normal for a data centre to be built from a plurality of separate data storage units to provide a data centre with expandable capacity and so that different parts of the data storage centre can be maintained independently. All, or part of, each of these data storage units may occasionally be offline due to a failure occurring or the data storage unit deliberately being taken offline for maintenance. The data is therefore preferably stored in such a way that it is still accessible when any of the data storage units for storing the data are offline. A known technique is to store the data in a replicated manner. That is to say, the same data stored in a data storage unit is also stored in one or more other, separate, data storage units. This approach allows fast retrieval of data from a back-up data storage unit if one of the primary data storage units is offline.\nA problem with storing data in a replicated manner is that the demand on data storage capacity is greatly increased and this increases costs. To solve this problem, it is known to use erasure codes for data storage. Using erasure codes for data storage requires additional processing to reconstruct unavailable data. However, the data storage requirements are a lot lower than those of replication systems. Reed-Solomon codes are well known error correcting codes that can be used as erasure codes for data storage. For a systematic implementation of a Reed-Solomon code as an erasure code, data is coded into a plurality of source data nodes, k, and a plurality of redundant data nodes, r, the total number of data nodes, n, being k+r. The source and redundant data nodes may be stored in a plurality of different data storage units, the data storage units each storing one or more data nodes. Due to the reconstruction properties of Reed-Solomon codes, if any r of the n data nodes are unavailable due to their data storage unit being offline, then the r data nodes can still be obtained through processing performed on a plurality of the other data nodes.\nA number of other coding techniques than Reed-Solomon are suitable for such data reconstructions. If the coding technique has the capability that any k nodes are able to reconstruct all of the source data nodes, the coding technique is Maximum Distance Separable, MDS. If the number of data nodes that are required to reconstruct a node is d, then for standard Reed-Solomon d=k.\nIt is highly preferable for the stored data to be coded systematically. This allows stored source data to be read directly from a data storage unit, without any processing being required to reconstruct the source data, if the data storage unit storing the source data is online.\nThere are known problems with using coded data for data storage as described above. For each data node that needs to be reconstructed, it is necessary to perform k read operations to obtain other data nodes and then processing operations to reconstruct the data nodes. The large number of read operations causes the data node reconstruction process to be slow and also results in a large increase in data traffic within the storage system.\nAn attempt to reduce the above-described problems of erasure codes is disclosed in the paper ‘A “Hitchhiker's” Guide to Fast and Efficient Data Reconstruction in Erasure-coded Data Centres’ by K. V. Rashimi et al, SIGCOMM 2014, Computer Communication Review, August 2014. In this paper a coding technique is disclosed that is motivated by a ‘piggy backing’ technique. The redundant data packets of two separate systematically coded sets of data packets are effectively combined with each other in a manner that reduces the necessary number of data reads. A first set of systematically coded data packets is generated using Reed-Solomon. A second set of systematically coded data packets is separately generated using Reed-Solomon. Each of the source data packets in the first set of data packets is respectively associated with source data packets in the second set of coded data packets and each of the redundant data packets in the first set of data packets is respectively associated with a redundant data packet in the second set of coded data packets. Each pair of associated data packets of the sets of coded data packets is stored together in the same data storage unit. The pairs of associated data packets are each stored in different data storage units. Accordingly, each data storage unit stores a pair of data packets, one from each set of coded data packets. The data packet form the first set of coded data packets is stored in a first sub-stripe of the data storage unit and the data packet form the second set of coded data packets is stored in a second sub-stripe of the same data storage unit, the first and second sub-stripes preferably not being contiguous. The paper provides results that demonstrate that the techniques for generating redundant data packets disclosed therein reduce the number of data read operations required from if the data had been coded using Reed-Solomon.\nOther techniques for coding data for storage in a data centre exist that all aim to reduce the amount of read data from what is necessary with standard Reed-Solomon coding. The coding techniques differ in how effective they are at improving on Reed-Solomon coding, the properties of the code that can be used (including the relative number of source and redundant data packets), the amount of processing required to generate the coded data and the amount of processing required to reconstruct source data from the coded data. For example, some coding techniques reduce the amount of data that needs to be read from that of Reed-Solomon but incur the drawback of it being necessary to store more data than if Reed-Solomon coding had been used. Other coding techniques do not allow the use of systematic codes and therefore incur the drawback of increasing the processing requirements.\nThere remains a need to improve on all known MDS coding techniques that have the advantageous properties of being systematic and not requiring a higher data storage capacity as standard Reed-Solomon."}
{"text": "The present invention relates to a multi-piece solid golf ball having a solid core composed of an inner core layer and an outer core layer, and having one or a plurality of cover layers encasing the core. More particularly, the invention relates to a multi-piece solid golf ball endowed with an excellent flight performance and an excellent spin performance.\nTo increase the distance traveled by a golf ball and improve the feel of the ball when played, efforts have hitherto been made to design golf balls with a multi-layer structure. Such efforts have led to the disclosure of various multi-piece golf balls in which the core, as well as the cover, has been provided with a two-layer structure. For example, JP-A 11-57070 and the corresponding JP No. 4006550 and U.S. Pat. No. 6,071,201 disclose a multi-piece solid golf ball having an inner core layer made of resin and an outer core layer made of rubber, wherein the inner core layer has a diameter of from 15 to 25 mm and a Shore D hardness of from 55 to 90, and the outer core layer has a JIS-C hardness of from 35 to 75 and a thickness of from 0.5 to 3.0 mm. However, because the inner core layer (center) of this golf ball is hard, the ball has a hard feel on impact and an increased spin rate on full shots.\nJP-A 2001-17571 and the corresponding U.S. Pat. No. 6,394,912 describe a golf ball in which the core is composed of a center core made of a thermoplastic resin or a thermoplastic elastomer and having a diameter of from 3 to 18 mm and a Shore D hardness of from 15 to 50, and of an outer core layer having a Shore D hardness near the boundary between the outer core layer and the center core that is from 1 to 15 units harder than the Shore D hardness of the center core. Also, JP-A 2000-229133 and the corresponding JP No. 3656806 and U.S. Pat. No. 6,605,009 disclose a golf ball having a construction composed of an inner core layer, an outer core layer and a cover, wherein the inner core layer is made primarily of resin and has a diameter of from 3 to 15 mm, the outer core layer is formed of a rubber composition, the center core has a Shore D surface hardness which is from 4 to 50 units harder than the innermost side of the outer core layer, and the specific gravities of these layers have been adjusted. However, this golf ball has a small center core and thus lacks a sufficient distance performance.\nU.S. Pat. No. 7,241,232 discloses a multi-piece solid golf ball having a multi-layer core in which an inner core layer is formed of a resin such as an ionomer, a polyamide or a polyester elastomer and the outer core layer is formed of a rubber composition, and having an inner cover layer and an outer cover layer composed of specific resins and having certain thicknesses. However, a sufficient distance is not achievable with this golf ball either.\nU.S. Pat. No. 7,468,006 discloses a golf ball having an inner core layer and an outer core layer, wherein the outer core layer is formed of a copolymer-based highly saturated ionomer having a Shore D hardness of 45 or more, and the inner core layer is formed of a terpolymer-based highly saturated ionomer having a Shore D hardness of 55 or less. However, in this golf ball, the outer core layer has been set to a higher hardness than the inner core layer and the ball does not have a high initial velocity. As a result, a sufficient distance is not achievable.\nJP-A 2008-301985 describes a golf ball having a ball construction of three or more layers in which a center core is made primarily of a thermoplastic resin and has a diameter of from 18 to 35 mm. However, the center core in this golf ball is soft, and so the ball does not have a high initial velocity. As a result, the distance on shots with a driver (W#1) leaves something to be desired."}
{"text": "The present invention concerns a method of producing long-length articles from hot-rolled carbon steel, such as, wire, strip and rolled sections.\nMethods of producing long-length articles from hot-rolled steel, such as high-strength wire, are well known in the prior art. For example, a hot-rolled wire rod is first subjected to heat treatment, whereupon the surface of the heat-treated rod is prepared for subsequent deformation by descaling and applying a sublubricant coat of lime, borax or copper. The rod is then subjected to cold plastic deformation, such as drawing, to obtain wire of a requisite cross-section.\nHeat-treating is employed for obtaining a metal structure which would enable subsequent cold deformation, accomplished in stages, to be effected with heavy overall reductions (up to 75-90%) which are needed to provide a preassigned strength.\nTo prevent wire embrittlement, cold plastic deformation is alternated with additional heat-treating operations (annealing or patenting).\nDepending on the designation of high-strength wire, the heat-treating operations (normalizing, patenting or quenching with subsequent tempering) as well as cold plastic deformation schedules, pass or overall reductions in particular, may vary during drawing.\nIn a number of cases the long-length articles -- wire -- are produced by only heat-treating the wire rod, e.g. by quenching with subsequent tempering, bainitizing, etc. However, wire made in the above manner, featuring perfect elongation and relaxation resistance, corrodes readily and has inadequate strength.\nThe wire produced by the use of only cold plastic deformation without preliminary heat treatment (such as patenting) possesses an adequate strength but has a very low state of plastic characteristics. Moreover, the drawing of such wire often leads to breaks because of its enhanced brittleness.\nKnown in the art is a method of producing long-length articles by plastic deformation of wire rods with simultaneous heating. Usually the above method is applicable to the so-called \"hard-to-work\" high-alloy (with chromium, titanium or tungsten) steels. However, in the case of using carbon steel as well, the metal should be always subjected to a preliminary heat-treating operation -- annealing or normalizing.\nWhen producing reinforcement high-strength wire, low-temperature tempering is used as a final operation to increase its ductility.\nAlso known is a method which consists essentially in that the long-length articles -- wire -- after preliminary heat-treating (patenting) and subsequent cold deformation are subjected to axial tension and concurrent heating to a temperature within the range of 300-400.degree.C. The latter operation, during which the processes of polygonization and blocking of dislocations by admixture atoms (nitrogen, oxygen, carbon) take place, contributes to the enhancement of rheological properties (relaxation resistance) of reinforcement and spring wire (see, e.g. British Pat. Nos. 748357, 969191, 969192. The wire subjected to the above operation features high physicomechanical characteristics and rheological properties.\nSimilar properties are obtainable by a more efficient method, wherein immediately after the last cold plastic deformation of a heat-treated (patented) wire rod, it is subjected to a subsequent plastic deformation combined with concurrent heating to a temperature within a range of 350.degree.-500.degree.C and axial tension.\nThis operation of the plastic deformation of wire heated concurrently to a temperature in the 300.degree.-600.degree.C range will be referred to as warm deformation.\nThus, all the known methods of producing high-strength wire are always characterized by heat-treatment of wire rod prior to drawing. The heat-treating operation is conducted on cumbersome and low efficiency equipment in heat-treating departments, and plastic deformation (drawing) is conducted in drawing shops. This fact creates considerable time loss and a high cost of in- or intrashop transportation resulting in a reduction in labour productivity and an increase in production cost."}
{"text": "1. Field of the Invention\nThis invention relates to a nitride semiconductor substrate and a method of fabricating the nitride semiconductor substrate. In particular, this invention relates to a nitride semiconductor substrate with a chamfered edge portion and a method of fabricating the nitride semiconductor substrate.\n2. Description of the Related Art\nIn a Si substrate, a GaAs substrate etc. as a semiconductor substrate used for fabricating an electronic device etc., the shape recognition and the localization of the substrate are conducted by image processing upon the transport during the device fabrication process, or upon the appearance inspection during the device fabrication process and the shipping. The shape recognition and the localization of the substrate are frequently conduced by irradiating the substrate with visible light or infrared light and detecting light reflected by the substrate.\nAs a conventional nitride semiconductor substrate, a gallium nitride (GaN) semiconductor substrate is known which is formed circular in its top view and has an edge portion with a surface roughness (Ra) of 10 nm to 5 μm (e.g., see JP-A-2004-319951).\nThe nitride semiconductor substrate as described in JP-A-2004-319951 is formed such that the edge portion is smoothened so as to reduce the incidence of cracks and to enhance the fabrication yield of an electronic device in the electronic device fabrication process using the nitride semiconductor substrate.\nHowever, since the nitride semiconductor substrate as described in JP-A-2004-319951 is transparent to (i.e., not reflective to) visible light and infrared light, it is impossible to achieve the shape recognition and the localization of the substrate by directly applying to the nitride semiconductor substrate the shape recognition method and the localization method used for a Si substrate, a GaAs substrate etc."}
{"text": "1. Field of the Invention\nThis invention relates to testing an electrical net of an electrical or electronic circuit and, more particularly, to a method of and apparatus for quickly detecting fault conditions in individual wiring nets in a multi-layer device carrier, for example. It also relates to a unique method and means for counting the number of pulses received from any source in a fixed period of time and to an interface which permits coupling to any general purpose computer for software manipulation of the collected test data.\n2. Description of the Prior Art\nIn the electrical arts, the testing of electrical \"nets\" of an electrical or electronic circuit for fault conditions is universally recognized to be of high importance. Among other things, circuit manufacturers and purchasers want to test nets for \"opens,\" \"shorts,\" and other faulty characteristics, such as circuits having incorrect resistance, inductance, and capacitance characteristics (\"RLC\" characteristics). The term \"net\" is known in the art and generally refers to the electrical characteristics of a circuit as they appear from a given circuit node. These electrical characteristics are measured with reference to a predetermined reference point, e.g., ground. A circuit typically includes several nets of interest.\nThe prior art includes methods that subject the net under test to a sinusoidal input stimulus and then measure both the magnitude and phase of the net's response to the input stimulus. The magnitude and phase may then be analyzed to detect faults in the net. In addition, the prior art includes methods that directly measure the various RLC characteristics of a net to be tested.\nThough the prior art methods are generally perceived as capable of detecting many types of faults, they have undesirable shortcomings. Among other things, certain nets, particularly those having large capacitances, require unduly long testing times. This is so, because the prior art methods must use a low frequency stimulus to ensure accurate measurement of a high capacitance. The low frequency results in longer testing times, and the longer testing times, in turn, increase the testing cost. This becomes particularly acute when high volume testing is required.\nIn addition, many prior art methods require a priori knowledge of the net configuration being tested, including knowledge of the arrangement of various components and their RLC characteristics. This complicates the testing, increases the cost, and in certain instances is entirely impracticable.\nFurther, many prior art methods are particularly configured to test nets having a relatively limited range of possible RLC characteristics. Consequently, manufacturers and purchasers are required to buy additional testing equipment, if they desire to test nets outside the limited range of their existing testers.\nIn addition, test equipment is typically designed for control by and operation with a specific computer for acquisition and evaluation of test data.\nAccordingly, there is a need in the art for an improved method and apparatus for testing nets quickly.\nIt is an object of the invention to provide a method and apparatus to detect faults in a net quickly.\nIt is a further object of the invention to provide a method and apparatus to detect faults in a net without requiring a priori knowledge of the net to be tested.\nIt is yet a further object of the invention to provide a method and apparatus that can test a variety of nets having a wide range of RLC characteristics.\nIt is a further object of the invention to provide a method and apparatus which can utilize any general purpose computer and its software for controlling the testing and fault detecting procedures thereof."}
{"text": "This invention relates to ballast apparatus for operating high-intensity-discharge (HID) high-pressure sodium vapor lamps and, more particularly, to a lead-type ballast having an improved power factor over the life of the lamp.\nThe conventional lead-type ballast has been utilized to successfully operate high-intensity-discharge lamps for many years. The circuit is characterized by a capacitor connected in series with the lamp. The capacitor functions to control lamp operation to provide inherent power factor correction. The lead-type ballast has been used, in recent years, to ballast HID sodium lamps. The HID sodium lamp is characterized by having a voltage that increases with age. As a result, limits in the maximum and minimum permissible lamp wattage and lamp voltage have been defined in the lamp industry as established by the American National Standards Institute (ANSI). These limits take the form of a trapezoid and HID sodium ballasts used today are designed to have the lamp operate within its boundaries over the life of the lamp.\nTo start the lamp, an electronic starting circuit is typically included as part of the ballast apparatus which provides pulses of sufficient voltage and duration until the lamp starts. The lead-type ballast is popular for use with the HID sodium lamp because it is intermediate in terms of cost and power losses while providing good wattage regulation. However, as the voltage of an HID sodium lamp increases with age the effect of the serially connected capacitor in maintaining power factor for the ballast apparatus is decreased.\nOther ballasts have been utilized to operate HID sodium lamps, such as the magnetic regulated or constant wattage ballast which includes a voltage regulating section which feeds a current limiting reactor which maintains the desired power factor through the majority of lamp life. The drawback of this ballast though is its relatively high cost and comparatively much higher power loss level than the loss level that can be attained using the lead-type ballast.\nOne such lead-type ballast is disclosed in German Pat. No. 1,802,011, dated Oct. 11, 1967, issued to Nuckolls. In the Nuckolls patent a high reactance transformer with a condenser is included in the secondary circuit, by which an adjustable degree of concentration is achieved in the secondary magnetic circuit that provides control of the lamp current to compensate for any changes in the supply voltage, as shown in FIG. 3 of the Nuckolls patent. A bypass condenser for high frequencies is also included across the primary of the higher reactance transformer."}
{"text": "Outriggers are popular in fishing applications. Outriggers are rigid poles that are a part of the boats rigging and are designed to extend the fishing pole bait lines beyond the “white wash” of the water surrounding the boat. This is preferable because bait traveling through clear water is more likely to successfully attract a fish. Once a fish takes the bait, the line is automatically released from the outrigger and all forces are transferred directly to a fishing pole secured to the gunnel of the boat. For example, the bait line can be connected to a quick release mechanism that releases the line and transfers the forces to the pole when a fish engages the bait or otherwise pulls on the line.\nOutriggers typically include several tubes connected to one another by clamps. The clamps can secure the tubes together when the outrigger is arranged in the extended position to prevent the tubes from collapsing upon one another. Also, the clamps can be engaged when arranged in a collapsed position to prevent the tubes from extending.\nOutriggers typically include clamps with guides that retain a filament. The filament can pass through the guides and connect to the quick release mechanism to extend the bait line a desired distance from the boat when the outrigger is in the extended position. The guides are preferably oriented parallel to one another. In doing so, the force applied to the filament by the bait line is coplanar with the axis of the clamp at a point furthermost from the boat, allowing a more structurally stable configuration.\nCurrent outrigger models require manual manipulation to retain the position of the clamp guides. For example, before engaging the clamps, a user must use their hands to rotate the tubes to properly align the guides of each clamp. In addition, conventional outrigger clamps and guides do not prevent the filament from becoming tangled, requiring the filament to be released and either untangled or replaced."}
{"text": "The present invention is in the technical field of firearms. More particularly, the present invention is in the technical field of ammunition magazine base plates. More specifically, the present invention is in the field of ammunition magazine base plates that can attach and detach more than one magazine.\nSince the inception of the guns, shooters have striven to improve the efficiency of re-loading. This was sensible for the hunter to avoid losing game, for the soldier to protect against his enemy, and for the farmer to defend his person and property. The revolver and repeating rifle improved upon the ball and black power by enabling the loading of multiple rounds, either in a rotating cylinder that could fire six bullets before re-loading was necessary, or into a captured enclosure in the rifle. The pistol improved upon the revolver by utilizing a magazine that can hold sometimes as many as 20 rounds, although 10 round magazines are by far the most common. For rifles one can find magazines that hold much larger numbers of rounds.\nToday, many states have instituted laws that limit the maximum number of rounds in a magazine to ten, and in some cases less than ten. The present invention helps individuals in regulated states as well as in unregulated states by securely fixing two magazines together at the base of the magazine by means of the present invention. This doubles the amount of ammunition attached and carried in a weapon. For some individuals, including professional and amateur competitive shooters, individuals protecting their families and property, military and law enforcement, the number of rounds in a magazine, and magazine replacement speed, can be critical.\nThe present invention provides for the attachment of two magazines of any capacity to a single base plate. A typical magazine includes the magazine body, a follower, a spring, an insert or locking plate and a base plate. The locking plate mechanism fixes into the base plate to prevent the magazine from disassembling while in use. The locking mechanism on a typical base plate cannot be released without easy access to the bottom of the magazine. Because of this design it has not previously been possible to connect two magazines at the base plates. The present invention employs a locking mechanism that can be operated without requiring direct access to the bottom of the magazines, and allows a user of the invention to benefit from being able to use the bullets from two magazines in a manner significantly more efficient that having to find a second magazine and insert it."}
{"text": "The invention relates to a process for the production of a vanadium/phosphorus mixed oxide catalyst precursor, its transformation into the active catalyst and a process for the production of maleic anhydride using this catalyst\nMaleic anhydride is a well known and versatile intermediate for manufacturing unsaturated polyester resins, pharmaceuticals or agrochemicals. It is usually produced by catalytic partial oxidation of aromatic (e. g., benzene) or non-aromatic (e. g., n-butane) hydrocarbons.\nThe main component of the active catalyst in the oxidation of non-aromatic hydrocarbons like n-butane to maleic anhydride is vanadyl pyrophosphate, (VO)2P2O7, which as a rule is obtained by thermal treatment of vanadyl acid orthophosphate hemihydrate of the formula (VO)HPO4xe2x80xa20.5H2O, acting as catalyst precursor. Both vanadyl pyrophosphate and vanadyl acid orthophosphate hemihydrate may, if desired, be accompanied by a promoter element selected from the groups IA, IB, IIA, IIB, IIIA, IIIIB, IVA, IVB, VA, VB, VIA, VIB and VIIIA of the periodic table of elements, or mixtures of such elements.\nMethods for preparing the precursor compound conventionally involve reducing a pentavalent vanadium compound under conditions which will provide vanadium in a tetravalent state (average oxidation number +IV).\nPrior art knows a great many different procedures, which however in general involve the use of vanadium pentoxide (V2O5) as a source of pentavalent vanadium and orthophosphoric acid (H3PO4) as the phosphorus source (see e.g. U.S. Pat. No. 5,137,860 or EP-A-0 804 963).\nAs a reducing agent in principle any inorganic or organic compound containing elements which are able to act as a redox couple possessing an oxidation potential suitable for this kind of reaction may be suitably applied.\nThe most common reducing agent is hydrogen chloride in aqueous solution.\nAlso favourably applied are organic media like primary or secondary aliphatic alcohols or aromatic alcohols such as benzyl alcohol as these compounds seem to at least in part dissolve the reactants and thereby facilitate the redox reaction.\nThe most preferred organic reducing agent is isobutyl alcohol as isobutyl alcohol combines optimal characteristics such as (i) a boiling point of 108xc2x0 C. at atmospheric pressure, (ii) dissolution of the vanadium alcoholates formed from V2O5, thus favouring a complete redox reaction in the liquid phase and (iii) achieving a redox potential for the couples isobutyl alcohol/isobutyraldehyde and isobutyl alcoholrisobutyric acid suitable to let the alcohol act as reducing agent. The tetravalent vanadium reacts with phosphoric acid (H3PO4) and leads to precipitation of the precursor vanadyl acid orthophosphate hernihydrate of the formula (VO)HPO4xe2x80xa20.5H2O. The precipitate is usually washed with isobutyl alcohol and subsequently dried.\nA major disadvantage of the conventional method as described above is that even after drying the precursor contains some percent of organic compounds from the organic reaction medium, compounds which are supposedly either (i) strongly adsorbed at the solid surface, and therefore not easily removable by the washing and drying treatment, or (ii) physically trapped in between the crystals of the precursor, or (iii) physically or chemically trapped a (xe2x80x9cintercalatedxe2x80x9d) in the crystalline structure of the precursor.\nIt has been found that this percentage of organic compound which remains trapped in the precursor is a fundamental parameter which can adversely affect the performance characteristics of the active catalyst obtained after the thermal treatment.\nThe object of the present invention therefore was to provide a method for controlling the carbon content in a vanadium/phosphorus mixed oxide catalyst precursor and accordingly to provide a superior catalyst precursor which, when activated, leads to superior results in the conversion of a non-aromatic hydrocarbon to maleic anhydride.\nIt was found that the objectives could be achieved with a new process for the preparation of a vanadium/phosphorus mixed oxide catalyst precursor according to claim 1.\nThe invention comprises reducing a source of vanadium in the presence of a phosphorus source in an organic medium which comprises\n(a) isobutyl alcohol or a mixture of isobutyl alcohol and benzyl alcohol and\n(b) a polyol\nin the weight ratio of 99:1 to 5:95.\nIn a mixture of isobutyl alcohol and benzyl alcohol the benzyl alcohol content is as a rule between 5 and 50 wt %.\nAs a source of vanadium a tetravalent or pentavalent vanadium compound may be applied. Representative examples, although not limiting, are vanadium tetrachloride (VCl4), vanadium oxytribromide (VOBr3), vanadium pentoxide (V2O3), vanadyl phosphate (VOPO4xe2x80xa2nH2O) and vanadium tetraoxide (V2O4). Vanadium pentoxide is the preferred vanadium source.\nAs mentioned above, the vanadium source may, if desired, be accompanied by promoter elements selected from the groups IA, IB, IIA, IIB, IIIA, IIIIB, IVA, IVB, VA, VB, VIA, VIB and VIIIA of the periodic table of elements, or mixtures thereof.\nPreferred promoter elements are selected from the group consisting of zirconium, bismuth, lithium, molybdenum, boron, zinc, titanium, iron and nickel.\nOrthophosphoric acid (H3PO4) is the preferred phosphorus source.\nIsobutyl alcohol is the preferred component (a).\nPolyols which can be used as the component (b) are expediently C2-6 aliphatic polyols, preferably C2-6-alkanediols such as 1,2-ethanediol, 1,2-propanediol, 1,3-propanediol, 1,2-butanediol, 1,3-butanediol, 1,4-butanediol, 2,3-butanediol, 1,2-pentanediol, 1,3-pentanediol, 1,4-pentanediol, 1,5-pentanediol, 2,3-pentanediol, 2,4-pentanediol, 1,2-hexanediol, 1,3-hexanediol, 1,4-hexanediol, 1,5-hexanediol, 1,6-hexanediol, 2,3-hexanediol, 2,4-hexanediol, 2,5-hexanediol and 3,4-hexanediol.\nMost preferred polyols are the C2-4-alkanediols 1,2-ethanediol, 1,3-propanediol and 1,4-butanediol.\nThe preferred mixture of alcohols contains 5 to 30 mol % of polyol with respect to isobutyl alcohol.\nAs a rule the vanadium source together with the phosphorus source is suspended in the organic medium and the mixture is kept under agitation at a temperature of expediently 90xc2x0 C. to 200xc2x0 C., preferably 100xc2x0 C. to 150xc2x0 C. over a period of 1 h to 24 h.\nThe ratio of vanadium source to phosphorus source is conveniently such that the P/V atomic ratio is in the range of 1:1 to 1.3:1, preferably 1.1:1 to 1.2:1.\nAs a rule after the reduction the precursor vanadyl acid orthophosphate hemihydrate of the formula (VO)HPO4.0.5H2O is formed which is filtered, washed and subsequently dried at a temperature of expediently 120xc2x0 C. to 200xc2x0 C.\nDue to the reduction treatment according to the invention the carbon content of the precursor can be controlled in the range of 0.7 wt. % to 15.0 wt. %, preferably in the range of 0.7 wt. % to 4 wt. %.\nIt has been found that best results are obtained with catalyst precursors which, after an additional thermal treatment at about 300xc2x0 C. for about 3 hours in air have a residual carbon content of 0.7 wt. % to 3 wt. %, most preferably 0.8 wt. % to 1.5 wt. %.\nOnce prepared the precursor can in view of its further activation treatment be formed into defined structures with defined properties. Such procedures may include wet grinding to a specific particle size, the addition of additives to improve attrition resistance, and the a formation of a convenient shape.\nA spherical shape for instance is most suitable for the application of the catalyst in a fluidized bed.\nThe further transformation of the so formed precursor into the active catalyst can be performed following a great number of activation processes known in the art, but in general include a heat treatment applying temperatures of up to 600xc2x0 C. More in detail, these processes may involve:\n(a) an initial heating of the precursor to a temperature not to exceed 250xc2x0 C.\n(b) a further heat treatment from about 200xc2x0 C. to at least 380xc2x0 C. to 600xc2x0 C. at the maximum\n(c) maintaining the temperature of stage (b) over a certain time and\n(d) cooling the activated catalyst, thereby maintaining an individually controlled atmosphere in all steps.\nIn a preferred embodiment, the activation of the precursor is accomplished using the procedure described in EP-A-0 804 963.\nIn a more preferred embodiment, the activation comprises the steps of\n(a) heating the catalyst precursor from room temperature to a precalcination temperature of about 300xc2x0 C. in air or oxygen-depleted air\n(b) keeping at said pracalcination temperature,\n(c) further heating the precalcined catalyst precursor in nitrogen up to a calcination temperature of about 550xc2x0 C. and\n(d) keeping at said calcination temperature.\nAfter the transformation into the active catalyst the catalyst is ready to be applied for the conversion of non-aromatic hydrocarbons to maleic anhydride.\nSuch processes are well known in the art, e. g. from U.S. Pat. No. 4,594,433, U.S. Pat. No. 5,137,860 or U.S. Pat. No. 4,668,652.\nIn general the non-aromatic hydrocarbon is converted with oxygen or an oxygen containing gas at a temperature from about 320xc2x0 C. to 500xc2x0 C. to maleic anhydride. The non-aromatic hydrocarbon is expediently selected from aliphatic C4-10 hydrocarbons, preferably nbutane. The conversion can take place in a fixed bed or a fluidized bed reactor but preferably in a fluidized bed reactor. The following examples are given by way of illustration only and arc not construed as to in any way limit the invention."}
{"text": "Integrated optical waveguides are often used as components in integrated optical circuits, which integrate multiple photonic functions. Integrated optical waveguides are used to confine and guide light from a first point on an integrated chip (IC) to a second point on the IC with minimal attenuation. Generally, integrated optical waveguides provide functionality for signals imposed on optical wavelengths in the visible spectrum."}
{"text": "The present invention relates to compounds with some structural similarities to bufadienolide compounds disclosed in the prior art. For a review of bufadienolides see Huimin Gao et al in Nat. Prod. Rep., 2011, 28, 953.\nThe bufadienolide compounds reported in the prior art are natural steroids, originally isolated from terrestrial natural sources such as plants of the families Crassulaceae, Hyacinthaceae, Iridaceae, Melianthaceae, Ranunculaceae, and Santalaceaethe, and animals of the genus Bufo (toads), Photinus (fireflies), and Rhabdophis (snakes) (Steyn et al. Nat. Prod. Rep. 1998, 15, 397-413; Krenn et al. Phytochemistry, 1998, 48(1), 1-29).\nAmong these bufadienolide compounds, Scilliroside and other Scilla compounds have been isolated from the red squill, Urginea maritima, and are disclosed to be highly toxic, specially Scilliroside which affects the cardiovascular and central nervous systems, causing convulsions and death (Verbiscar et al. J. Agric. Food Chem. 1986, 34, 973-979; Kopp et al. Phytochemistry, 1996, 42(2), 513-522). Majinda et al. also isolated bufadienolide compounds from Urginea sanguinea, which make the plant unsafe to be used as a medicinal plant (Planta Med. 1997, 63, 188-190).\nThe antiviral activity against a series of rhinoviruses and the anti-herpetic activity of several bufadienolide compounds were evaluated by Kamano et al. (Chem. Pharm. Bull. 1988, 36(1), 326-332) and Takechi et al. (Phytochemistry, 1996, 41(1), 125-127), respectively, and they found that most of the compounds displayed some inhibitory activity.\nAdditionally, the cytotoxic activity of several bufadienolide compounds has been evaluated by several authors. Specifically, Jing et al. reported that bufalin has a potent growth-inhibitory effect on human cancer cells of leukaemia (HL-60, ML1, U937, and K562 cell lines), epithelioid carcinoma (HeLa cell line), hepatoma (PLC/PRF/5 cell line), and epidermoid carcinoma (A431 cell line), but it is less potent on mouse leukaemia M1, melanoma B16, and lymphoid neoplasm P388 cell lines and rat hepatoma AH66 and chromaffin cell PC12 cell lines. They also found that bufalin induces typical apoptosis in human leukaemia HL-60 cell line but not in human leukocytes (Jpn. J. Cancer Res. 1994, 85(6), 645-651).\nKupchan et al. disclosed several bufadienolide compounds isolated from Bersama abyssinica which showed inhibitory activity against human carcinoma of the nasopharynx (KB) cell line (Bioorg. Chem. 1971, 1, 13-31; J. Org. Chem. 1971, 36(18), 2611-2616).\nKamano et al. evaluated the cytotoxic activity of 80 bufadienolide and cardenolide compounds, isolated from the Chinese drug Ch'an Su (obtained from the skin glands of toads such as Bufo gargarizans), against a primary liver carcinoma PLC/PRF/5 cell line and the colchicine-resistant cell line of PLC/PRF/5. Of them, 16 were shown to have potent cytotoxicities (IC50<10−3 μg/mL) against PLC/PRF/5 cell line (Bioorg. Med. Chem. 1998, 6, 1103-1115; J. Med. Chem. 2002, 45, 5440-5447). Additional bufadienolide compounds were isolated by Nogawa et al. from the same source, which were tested against human carcinoma of the nasopharyx (KB), human leukaemia (HL-60), murine leukaemia (MH60), pancreas adenocarcinoma (BXPC3), breast adenocarcinoma (MCF7), CNS glioblastoma (SF268), lung NSC (NCIH460), colon carcinoma (KM20L2), and prostate cancer (DU145) cell lines (J. Nat. Prod. 2001, 64, 1148-1152).\nYe et al. prepared novel bufadienolide compounds from bufalin by microbial hydroxylation. The compounds were tested against human hepatoma Bel-7402, human gastric cancer BGC-823, human cervical carcinoma HeLa, and human leukaemia HL-60 cell lines, showing some of them potent cytotoxicities comparable to that of bufalin (J. Steroid. Biochem. Mol. Biol. 2004, 91, 87-98; J. Nat. Prod. 2005, 68, 626-628).\nSince cancer is a leading cause of death in animals and humans, several efforts have been and are still being undertaken in order to obtain an anticancer therapy active and safe to be administered to patients suffering from a cancer. The problem to be solved by the present invention is to provide compounds that are useful in the treatment of cancer."}
{"text": "This invention relates to coated polyester fibers which are indigo-dyeable and possess wash-down and crocking resistance characteristics similar to indigo-dyed cotton fibers and to a process for preparing such fibers.\nIndigo-dyed denim fabrics are twill fabrics in which only the warp yarns are dyed. For reasons of style, it is desired not only that these fabrics have an initial deep indigo blue color, but also that the fabrics begin to fade in an obvious manner after only a few wearings and launderings. For stronger, more durable fabrics with better fabric stability than all-cotton fabrics, blends of polyester staple fibers with cotton have been used. The undyed filling yarns can be made of 50/50 polyester/cotton blends for high strength. However, since commercially available polyester fibers will not dye with indigo, it has been found that the warp yarns can contain only a small amount of polyester fibers--only about 10% wt. % in open-end-spun yarns and no more than about 25wt. % in ring-spun yarns--if the desired deep blue color is to be obtained. At higher blend levels, it becomes increasingly difficult to make a fabric which can be dyed to an acceptable indigo shade.\nThe use of polyaminoalkylsilanes to coat either natural or synthetic fibers so that they will be dyeable with a wide variety of dyestuffs has been disclosed by Speier in his U.S. Pat. No. 3,504,998. The polymerization of unsaturated compounds containing polyalkylene oxide segments onto polyester, polyamide, and polyacrylonitrile fibers to provide them with indigo-dyeable surfaces is described by Toray Industries, Inc., in their Japanese Patent Application (Kokai) 77778/79. However, polyester staple fibers adapted for indigo dyeability with wash-down and fading characteristics adequately simulating indigo-dyed cotton fibers have not been described by the prior art."}
{"text": "The present invention relates generally to techniques which are effectively applied for improving the characteristics of a voltage controlled oscillation circuit (VCO) capable of switching from one oscillating frequency to another as well as the characteristics of an on-chip VCO, and facilitating measurements of the characteristics of such VCOs, and more particularly, to techniques which are effectively utilized in a VCO mounted in a high frequency semiconductor integrated circuit for demodulating a reception signal and modulating a transmission signal in radio communication apparatuses, for example, a portable telephone and the like which can transmit and receive signals in a plurality of bands.\nA radio communication system such as a portable telephone uses a PLL (phase locked loop) circuit which has a VCO for generating an oscillating signal at a predetermined frequency. The oscillating signal is combined with a reception signal and a transmission signal. Conventional portable telephones include a dual-band portable telephone which can handle signals in two frequency bands, for example, a GSM (Global System for Mobile Communication) signal in a band of 880-915 MHz and a DCS (Digital Cellular System) signal in a band of 1710-1785 MHz. Some dual-band portable telephones are designed to support two different bands with a single PLL circuit by switching the frequency of the PLL circuit.\nIn recent years, however, a need exists for a triple-band portable telephone which can handle, for example, a PCS (Personal Communication System) signal in a band of 1850-1915 MHz in addition to the GSM and DCS signals. It is also contemplated that the portable telephones are required to support a larger number of bands in the future.\nFor a high frequency semiconductor integrated circuit (hereinafter called the “high frequency IC”) designed to modulate a transmission signal and demodulate a reception signal, for use in such a portable telephone which can support a plurality of bands, a direct conversion system is effective from a viewpoint of a reduction in the number of parts. While the direct conversion system is relatively easy in supporting a plurality of bands, a VCO should be capable of oscillating over a wide frequency range. In this event, when a single VCO is used with the intention to cover the overall frequency range, the resulting VCO would be extremely sensitive to a control voltage applied thereto, and therefore vulnerable to extraneous noise and fluctuations in a power supply voltage.\nOn the other hand, a reduction in the number of parts may be effectively accomplished by forming a VCO, which has been typically fabricated in a module separate from a high frequency IC in many cases, on the same semiconductor chip on which the high frequency IC is fabricated. However, since an on-chip VCO manufactured by the current technologies experiences large variations in the absolute value of the oscillating frequency, the on-chip VCO must be provided with a function of correcting the oscillating frequency after the manufacturing. However, if the variations are corrected by trimming based on a mask option or a bonding wire option, typically used in conventional semiconductor integrated circuits, the cost is inevitably increased."}
{"text": "Psychosis such as schizophrenia is a severe and crippling mental disorder that affects about 1% of the population. It is a mental disorder that is characterized by gross impairment in reality, major disturbances in reasoning, often evidenced by delusions and hallucinations, incoherent speech, and/or disorganized and agitated behavior. Several classes of anti-psychotic drugs are available for treatment of schizophrenia, including the prototypical antipsychotic drugs such as chlorpromazine and haloperidol as well as many others such as droperidol, fluphenazine, loxapine, mesoridazine molidone, perphenazine, pimozide, prochlorperazine promazine, thioridazine, thiothixene, and trifluoperazine. While these agents are effective in treating positive symptoms of psychosis such as symptoms of hallucination and delusions, e.g., in schizophrenia, these drugs often cause both short-term and long-term movement disorders and other side effects including acute dystonia (e.g., facial grimacing, torticollis, oculogyric crisis, abnormal contraction of spinal muscles and of muscles involved in breathing), akathisia, bradykinesia, rigidity or short term paralysis, parkinsonism, sedation, dry mouth, sexual dysfunction and sometimes tardive dyskinesia. Tardive dyskinesia may persist after discontinuation of the use of typical antipsychotic agents and there is no effective treatment of such side effects. Because of the severity of the side effects, typical antipsychotic drugs, though effective in treating the mental and emotional aspect of the disorder, do not help patients to function normally in society.\nAlthough another class of antipsychotic agents called atypical antipsychotic agents, which include clozapine, aripiparazole, olanzapine, quetiapine, risperidone and ziprasidone (atypical antipsychotic agents) are effective in treating positive and negative symptoms of schizophrenia with fewer extrapyramidal side effects, these agents can nevertheless cause other serious and at times fatal side effects, including bone marrow suppression, seizure, orthostatic hypotension, insomnia, sedation, somnolence, weight gain, and if administered at higher doses, may again cause extrapyramidal side effects. Therefore, atypical antipsychotic agents, though have improved clinical profiles, are nevertheless undesirable.\nIn addition to the positive and negative symptoms of psychosis (e.g., schizophrenia), many psychotic patients often times also suffer from depression. While both typical and atypical antipsychotic agents are effective in treating psychosis, depression is often times neglected or left under-treated. The combination of psychosis and depression poses a particular challenge in their treatment as studies revealed that up to 10% of the patients suffering from schizophrenia end their own lives. Therefore, there is a need for agents that are useful for the treatment of psychosis in depressed patients, and for the treatment of depression as well as other disorders such as sleep and mood disorders in psychotic patients and patients suffering from Parkinson's disease without exhibiting or exhibiting minimal extrapyramidal and other side effects compared to conventional antipsychotic, hypnotic and anti-depressive agents.\nSubstituted heterocycle fused gamma-carbolines are known to be agonists or antagonists of 5-HT2 receptors, particularly 5-HT2A and 5-HT2C receptors, in treating central nervous system disorders. These compounds have been disclosed in U.S. Pat. Nos. 6,548,493; 7,238,690; 6,552,017; 6,713,471; U.S. RE39680, and U.S. RE39679, as novel compounds useful for the treatment of disorders associated with 5-HT2A receptor modulation such as obesity, anxiety, depression, psychosis, schizophrenia, sleep disorders, sexual disorders migraine, conditions associated with cephalic pain, social phobias, and gastrointestinal disorders such as dysfunction of the gastrointestinal tract motility. PCT/US08/03340 and U.S. application Ser. No. 10/786,935 also disclose methods of making substituted heterocycle fused gamma-carbolines and uses of these gamma-carbolines as serotonin agonists and antagonists useful for the control and prevention of central nervous system disorders such as addictive behavior and sleep disorders. Although these references disclose use of substituted heterocycle fused gamma-carbolines to independently treat disorders associated with serotonin pathways such as sleep disorder, depression, psychosis, and schizophrenia associated with the 5-HT2A pathways, there is no teaching that specific compounds of substituted heterocycle fused gamma-carbolines also exhibit nanomolar binding affinity to serotonin reuptake transporter (SERT) and dopamine D2 receptors and therefore may be used to treat a combination of psychosis and depressive disorders as well as sleep, depressive and/or mood disorders in patients with psychosis or Parkinson's disease.\nIn addition to disorders associated with psychosis and/or depression, these references do not disclose use of particular substituted heterocycle fused gamma-carbolines at a low dose to selectively antagonize 5-HT2A receptors without affecting or minimally affecting dopamine D2 receptors, thereby useful for the treatment of sleep disorders without the side effects of the dopamine D2 pathways or side effects of other pathways (e.g., GABAA receptors) associated with convention sedative-hypnotic agents (e.g., benzodiazepines) including but not limited to the development of drug dependency, muscle hypotonia, weakness, headache, blurred vision, vertigo, nausea, vomiting, epigastric distress, diarrhea, joint pains, and chest pains."}
{"text": "This invention relates to lightweight anchors. More particularly, this invention relates to a device for mounting a lightweight anchor on the rail of a boat.\nA lightweight anchor is a type of marine anchor that is well known in the art and widely used with boats of all sizes. The component parts of a lightweight anchor include a stock, a crown, a pair of flukes which are fixed relative to the crown and a shank which is mounted for pivotal movement relative to the flukes. In my copending patent application Ser. No. 746,100, filed concurrently with this patent application, the subject matter of which patent application is incorporated herein by reference, there is described one type of lightweight anchor construction. Other examples of lightweight anchors can be found in U.S. Pat. Nos. 3,780,688; 3,782,318; and 3,783,815.\nWhen not being used, lightweight anchors are usually either stashed away underneath the deck, secured to the top of the deck or attached to some part of the boat, such as the rail. Hitherto, lightweight anchors have been mounted on the rail using a pair of metal straps, one of the straps having a hole and the other strap having a slot. The straps are secured to the rail and the anchor is mounted on the straps by inserting one end of the stock into the hole and then dropping the stock on the other side of the flukes into the slot. One of the disadvantages of this type of mount is that the stock is free to rotate within the hole and slot about its own longitudinal axis. This results in an undesirable swaying or pendulum type movement of the flukes, even under mild sea conditions. To prevent this rotation the flukes must be lashed down. Another disadvantage of this type of mount is that the shank is also free to swing back and forth in a pendumum type movement since it is pivotally mounted relative to the stock. To prevent such rotation, the shank must be tied down."}
{"text": "1. Field of the Invention\nThis invention concerns a device for the continuous supply of liquid mediums having varying viscosities, in particular of lubricants.\n2. Description of the Related Art\nA grease box is known with a closed expandable body and a space filled with lubricant whereby the sealed expandable body contains a device which can be triggered at any monment and consists of a galvanic element for producing gas which builds up pressure in the expandable body where it expands by a certain amount every day and, in doing so, dispenses lubricant via a lubricant channel from a lubricant reservoir to a place in need of lubrication (DE-PS 1256001). It is possible that in the known grease box--if the expandable body is not gas-proof--part of the compresed gas will escape and find its way into the grease container. It is known that this can be avoided by installing a separation device between the expandable body and the lubrication space, which separation device slides along the cylindrical inside wall of the grease box and seals it (DE-PS 2139771).\nA common element in these known embodiments is that the sealed expandable body contains an electrolytic liquid in which a corrosive element is inserted by means of a screw to trigger the development of gas. Apart from the fact that the electrolytic liquid requires a large space and involves the danger of leakage and therefore damage to the environment, once the formation of gas has started in these known embodiments, it cannot be stopped or interrupted. It is not possible to replace the electrolyte and/or lubrication containers without destroying the grease box. Similarly, it is not possible to simply add more lubricant to the container.\nThe purpose of this invention is to simplify the basic design of the known devices, to enlarge the space for the medium to be dispensed in prportion to the gas-generating unit, to make the operation easier, to simplify the regulation of the amount and time involved in the generation of gas, and to allow for easy retrofitting while, at the same time, avoiding the danger of electrolyte leakage.\nThe use of a compact cell reduces considerably the space requirements for the gas-generating part and eliminates the danger of electrolyte leakage. The insertable component allows not only for the possibility of replacement when the cell is exhausted for the addition of the medium to be dispensed, e.g., the lubricant, but the component also functions, in conjunction with corresponding elements of its counterpart, as a \"switch\" for starting and stopping the generation of gas without, as a rule, requiring a special switch.\nFor this purpose the cell is mounted as a moving part in the component and is inserted into the electric circuit by means of, for example, a projection on the container lid against which the cell comes to rest when the component is screwed into place. Consequently, it does not require the introduction of plugs or similar to close the electric circuit.\nThe secondary claims describe additional measures for the effective execution and further elaboration of the device according to the invention.\nFor example, the purpose of the measure according to the present invention is to prevent the unintended shifting of the cell in the component. This may involve a spring, or a (spring-loaded) projection with comparable action.\nThe measure according to the present invention for example, creates in a simple manner a device which triggers the generation of gas when the component is inserted. This may also be achieved--possibly as an additional measure--by means of a switch for the closing of the electric circuit.\nTo prevent the compressed gas from escaping during a temporary stoppage or interruption of the gas generation, a check valve is provided which also serves to protect the device during transportion until the moment of installation, i.e., in its condition for delivery. Because the gas yield of the compact cell is limited it is possible, according to the present invention to keep part of the water needed for the gas generation outside the cell in a separate space. Accordingly hereinafter, this space can be filled with a water-absorbent mass. Due to the difference in vapor pressure between the cell and the mass--which has preferably an atmospheric pressure which is low but similar to that of the electrolyte--water from the cell can be supplied as needed from the cell to the mass or, in reverse, from the mass to the cell. In addition, the mass offers the advantage of filling a hollow space, thereby reducing the amount of atmospheric oxygen commonly present in hollow spaces. Consequently the cell--which normally when started up uses first the atmospheric oxygen before generating gas--will switch much sooner to its specific purpose, i.e., the generation of gas. To prevent electric transition resistance, additional material may be added to the mass.\nThe same purpose of increasing the gas yield is also served by the measure described hereinafter, according to which the device is equipped with two or more cells which are insulated from each other upon delivery but are comnnected electrically with one another when the component is inserted into the container lid. For this purpose a switch may be mounted between the cells and/or the switch itself may be designed, according to the present invention, as a flexible spring element separating and insulating the cells electrically before insertion, while allowing them to be electrically connected after insertion. Obviously, the cells may also be switched on or connected to the electric circuit in sequence, i.e., the second cell is activated only when the first one is used up.\nAdditional especially designed embodiments of the device according to the invention, or individual components are described hereinafter. Of course, the devices, measures and method of connection described therein can be substituted by equivalent devices, measures or methods of connection.\nBesides the aforementioned characteristics/advantages of the solution according to the invention, the following is worth mentioning:\nThe invention is not limited to its use as a grease box, it is also capable of dispensing viscous mediums of another nature. PA1 According to the requirements, cells of various types and capacity may be used. There is also the possibility, as mentioned, of using several cells in one component, which cells may operate simultaneously or in sequence. PA1 As long as the component has not been screwed into place, no gas will be formed. To avoid unexpected triggering during storage or transportation, the bore holding the cell can be closed off with a small cover or similar. PA1 Before delivery, the bore in the container lid may be closed off with a fastening closure which is removed when the unit is put into operation. PA1 To increase and/or accelerate the generation of gas, an additional external power source may be provided."}
{"text": "Tremolo bridges having tremolo activation levers, which are often called whammy bars, vibrato bars, or tremolo bars, that are designed so that the musician can easily add vibrato while playing the guitar. Vibrato is the periodic variation in a tone and tremolo is periodic variation in the amplitude. Because of this confusion, tremolo and vibrato will be used interchangeably herein. This is accomplished by having the base plate of the bridge pivoting on a fixed point or points. The vibrato arm is connected to the base plate and when the musician moves the vibrato arm the strings of the guitar (or other instrument) increase and decrease in tension. As a result, the tone varies periodically resulting in a vibrato effect.\nOne problem with tremolo bridges is that when a musician pulls (bends) on a string, the tension of the other strings are affected. This occurs because the base plate is balanced on pivot point (line or anchor) and increasing the tension on one string results in base plate moving from its resting position. When the base plate moves the tension on the other strings change. In most cases it causes the tuning of the other strings to go slightly flat. This is an undesirable side effect of tremolo bridges and negatively effects the quality of musical sound.\nThere have been attempts to solve this problem by adding counterbalancing springs with different spring constants to the tremolo bridge. However, these have been ineffective since the tremolo bridge is already balanced between the force of the strings and the springs on the tremolo bridge. Adding additional springs does not solve the problem.\nThus there exists a need for a tremolo bridge that does not result in the other strings going flat when one of the strings is tensioned. Note that while the present invention is described with respect to guitars, it is applicable to any string instrument."}
{"text": "RMI is a set of protocols developed by Sun Microsystems, Inc. of Palo Alto, Calif. which enable Java objects (based on the Java programming language developed by Sun Microsystems, Inc.) to communicate remotely. The use of RMI allows for distributed computing in which the processing being performed by an application may be divided between a client and one or more servers. The RMI services being invoked are managed by an RMI activation system daemon such as RMID or Phoenix. RMID is the RMI activation system daemon from Sun Microsystems, Inc. (started by the ‘rmid’ command) that allows objects to be registered and activated in a Java Virtual Machine (JVM). Phoenix is a distributed network framework developed by Intel Corporation of Santa Clara, Calif. that is backwards compatible with RMID and may be used to support the RMI service framework discussed herein. The RMI activation system daemon provides persistence to the RMI services by managing the process in which the RMI service runs. In the event of session failure, the RMI activation system daemon is able to reinitiate the RMI service to limit any interruption in service being provided to the client.\nUnfortunately, the RMI activation system daemon uses a number of mechanisms which are Java-based. While this allows the RMI services being managed by the RMI activation system daemon to communicate with each other using well-understood Java mechanisms, the use of the Java mechanisms presents a problem in integration with non-Java applications. Conventionally, there is not an available mechanism by which non-Java applications may be integrated into a service framework such that it can be managed by the RMI activation system daemon."}
{"text": "There is an emerging need to provide manufactured products of all types in which the levels of toxins is minimized or completely eliminated. This need has a clear underlying medical basis and is accelerated by fear and pending legislation—the results of many recent well-publicized cases of toxins in manufactured products (e.g., lead in toys). The costs of unsafe products go well beyond the health impact to include significant loss of business, permanent damage to brands and corporate image, and increased levels of corporate and personal liability.\nIn response to these problems, there is a growing trend of increasingly strict environmental and health regulations of consumer products around the world. The list of products regulated is rapidly increasing and the types and permitted levels of toxins are becoming more restrictive. Some industry players are going beyond the regulations for the products they distribute, by mandating even cleaner products in their supply chain. Regulations are effectively aimed at decreasing direct human exposure to toxins by reducing toxins in our environment. Several stricter standards can be traced to European environmental directives that began in the early 1990s, starting with regulations in packaging materials and batteries. In subsequent years, reductions on hazardous substances were introduced by the EU for automobiles (ELV) and two directives related to electronics (Restriction of Hazardous Substances or RoHS and Waste Electrical and Electronic Equipment or WEEE). Pending U.S. Federal legislation lowers allowable lead levels in paint on toys by a factor of six and threatens criminal prosecution for companies that violate with penalties ranging from $10 million to $100 million for a single violation. In addition, nine other known toxins are targeted for restriction, including: mercury, arsenic, cadmium, barium, and chromium.\nThe spread of such human health and environmental initiatives are having profound global implications on the way products are designed, manufactured, and ultimately discarded or recycled.\nCurrent measurement methods for toxins in manufactured products do not meet the needs of the supply chain, from the factories to the ultimate consumers. Identification and measurement of toxins are needed at each step of the chain, from raw materials, to components, to finished goods. While raw-material measurements are most efficient for factories, distribution channels typically require measurements on the final product. New techniques are urgently needed to accurately, quickly, consistently, and cost-effectively measure toxins at each stage, with minimal interruptions in the flow of manufacturing and distribution of the goods. Because toys and other manufactured products often have small painted features (pigments are often the source of the toxins), it is necessary to measure small areas while differentiating the paint from the base material.\nExisting low-cost methods of toxin detection are generally ineffective, e.g., swab tests. Higher-cost methods that provide the requisite accuracy are expensive and time consuming. These sometimes involve: manually scraping samples, digesting them in acids at elevated temperature and pressure, introducing them into a combustion chamber, and analyzing the combustion product. One widely used method today is inductively coupled plasma optical emission spectroscopy (ICP-OES)—a method which is expensive, destructive, and slow. Alternatively, handheld x-ray fluorescence (XRF) guns are rapid and nondestructive, but are only reliable for higher than regulated concentrations, and are averaged across large sample areas, and cannot separately evaluate paint layers.\nAs discussed further below, the present invention provides a measurement solution having fast, accurate results for toxins in manufactured products, enabled by sophisticated proprietary x-ray optics. Such proprietary optics typically provide 10-1,000× improvements in the ability to focus x-rays; and optic-enabled analyzers are especially suited for these targeted markets—moving measurements from the lab into the factory, field, and clinic.\nIn x-ray analysis systems, high x-ray beam intensity and small beam spot sizes are important to reduce sample exposure times, increase spatial resolution, and consequently, improve the signal-to-background ratio and overall quality of x-ray analysis measurements. In the past, expensive and powerful x-ray sources, such as rotating anode x-ray tubes or synchrotrons, were the only options available to produce high-intensity x-ray beams, in the laboratory. Recently, the development of x-ray optic devices has made it possible to collect the diverging radiation from an x-ray source by focusing the x-rays. A combination of x-ray focusing optics and small, low-power x-ray sources can produce x-ray beams with intensities comparable to those achieved with more expensive devices. As a result, systems based on a combination of small, inexpensive x-ray sources, excitation optics, and collection optics have greatly expanded the availability and capabilities of x-ray analysis equipment in, for example, small laboratories and in the field, factory, or clinic, etc.\nMonochromatization of x-ray beams in the excitation and/or detection paths is also useful to excite and/or detect very precise portions of the x-ray energy spectrum corresponding to various elements of interest (lead, etc.). X-ray monochromatization technology is based on diffraction of x-rays on optical crystals, for example, germanium (Ge) or silicon (Si) crystals. Curved crystals can provide deflection of diverging radiation from an x-ray source onto a target, as well as providing monochromatization of photons reaching the target. Two common types of curved crystals are known as singly-curved crystals and doubly-curved crystals (DCCs). Using what is known in the art as Rowland circle geometry, singly-curved crystals provide focusing in two dimensions, leaving x-ray radiation unfocused in the third or orthogonal plane. Doubly-curved crystals provide focusing of x-rays from the source to a point target in all three dimensions. This three-dimensional focusing is referred to in the art as “point-to-point” focusing.\nCommonly-assigned U.S. Pat. Nos. 6,285,506 and 7,035,374 disclose various configurations of curved x-ray optics for x-ray focusing and monochromatization. In general, these patents disclose a flexible layer of crystalline material (e.g., Si) formed into curved optic elements. The monochromating function, and the transmission efficiency of the optic are determined by the crystal structure of the optic.\nThe ability to focus x-ray radiation to smaller spots with higher intensities, using focusing and monochromating x-ray optics, has enabled reductions in the size and cost of x-ray tubes, and x-ray systems have therefore been proliferating beyond the laboratory to in-situ, field uses. Commonly-assigned U.S. Pat. Nos. 6,934,359 and 7,072,439, incorporated by reference herein in their entirety, disclose monochromatic wavelength dispersive x-ray fluorescence (MWD XRF) techniques and systems, using doubly curved crystal optics in the excitation and/or detection paths. The x-ray optic-enabled systems described in these patents have enjoyed widespread success beyond the laboratory, for measuring sulfur in petroleum fuels in a variety of refinery, terminal, and pipeline environments.\nIn such systems, precise optic alignment along an axis defined by a source and sample spot may be required, as illustrated in above-incorporated U.S. Pat. No. 7,035,374, which proposes an arrangement of curved, monochromating optics around a central axis operating according to Bragg diffraction conditions. FIG. 1a is a representative isometric view of this x-ray optic arrangement 150 having a curved optic 152, an x-ray source location 154, and an x-ray target location 156. X-ray source location 154 and x-ray target location 156 define a source-to-target transmission axis 162. Optic 152 may include a plurality of individual optic crystals 164, all of which may be arranged symmetrically about axis 162.\nFIG. 1b is a cross-sectional view taken along section lines 1b-1b of FIG. 1a, wherein the surface of optic 152, x-ray source location 154, and x-ray target location 156 define one or more Rowland (or focal) circles 160 and 161 of radius R for optic 152. Those skilled in the art will recognize that the number and orientation of the Rowland circles associated with crystal optic 152, or individual crystals 164, will vary with the position of the surface of optic crystal 152, for example, the variation of the toroidal position on optic crystal 152.\nThe internal atomic diffraction planes of optic crystal 152 also may not be parallel to its surface. For example, as shown in FIG. 1b, the atomic diffraction planes of crystal 152 make an angle γ1 with the surface upon which x-rays are directed, at the point of tangency 158 of the surface and its corresponding optic circle 160 or 161. θB is the Bragg angle for crystal optic 152 which determines its diffractive effect. Each individual optic crystal can in one example be fabricated according to the method disclosed in above-incorporated U.S. Pat. No. 6,285,506, entitled “Curved Optical Device and Method of Fabrication.”\nAll individual crystals 164 should be aligned to the source-to-target axis 162, for proper Bragg conditions. Improvement in optic alignment, especially for such multiple-crystal optics, therefore remains an important area of interest. Another issue, which particularly affects volume manufacturing, is the need to align disparate components which may be purchased from different vendors. For example, the x-ray tubes, when purchased in quantities from a vendor, may have source x-ray spots which are not consistently centered relative to their own housings. Re-centering these x-ray tube spots is necessary, as an initial step in the alignment process for an entire x-ray source assembly.\nVarious optic/source combinations have already been proposed to handle thermal stability, beam stability, and alignment issues, such as those disclosed in commonly assigned U.S. Pat. Nos. 7,110,506; 7,209,545; and 7,257,193. Each of these patents is also incorporated herein by reference in its entirety. In particular, U.S. Pat. Nos. 7,209,545 (entitled “X-Ray Source Assembly Having Enhanced Output Stability, and Fluid Stream Analysis Applications Thereof”) and 7,257,193 (entitled “X-Ray Source Assembly Having Enhanced Output Stability Using Tube Power Adjustments and Remote Calibration”) address certain tube/optic alignment problems during source operation with real-time, corrective feedback approaches for alignment between the tube focal spot, optic, and output focal spot. Sensors are used to detect various operating conditions, and mechanical and/or thermal adjustments are made to correct for instabilities, including misalignments. These types of systems are necessary and valuable for certain applications, but can also increase the cost and complexity of fielded systems.\nThe above-described XRF technology and systems have been useful in single element analyzers for measuring generally homogeneous sample structure (e.g., sulfur in petroleum products). However, the measurement of toxins in manufactured products presents an additional level of challenges. First, an instrument should have the capability to measure more than one element simultaneously or near-simultaneously, from a relatively confined list of about 10 toxic elements, discussed above. Moreover, manufactured products are likely to be heterogeneous in nature, requiring small spot resolution, as well as the ability to detect toxins in one of a number of heterogeneous layers (e.g., the level of lead in a paint layer and a substrate layer beneath the paint).\nImproved x-ray analysis method and systems are required, therefore, to address the problems associated with measuring multiple toxins in potentially heterogeneous samples, to enable in-the-factory and/or in-the-field measurement of toxins in manufactured products."}
{"text": "1. Field of the Invention\nThe present invention relates to a therapeutic device for improving the respiration of a patient, with a curved or bent pipe section and a mouthpiece inserted in its first end.\n2. Description of the Prior Art\nA therapeutic device for supporting the respiration of a patient is disclosed in EP 0681853. In this case, the curved pipe section has a height-adjustable mouthpiece inserted in its end. As a result, the mouthpiece penetrates inside the pipe section in areas. The penetrating area of the mouthpiece has a hose attached to it, which is bent by the inner contour of the pipe section. The second end of the pipe section is open, therefore the therapeutic device can be used for inhalation and exhalation. The ending of the hose section causes the device to undergo oscillatory vibration during inhalation and exhalation, thereby causing vibrations to be generated in the pharyngeal and pulmonary areas of the patient, thus allowing diseases of the airways to be treated.\nAlthough the therapeutic device disclosed in EP 0681853 A1 has proven effective in practice, it has become apparent that no optimum treatment successes can be achieved when the air is inhaled through the hose section because the vibrations generated by the hose section during inhalation are significantly smaller in amplitude and frequency than those during exhalation.\nA further disadvantage with this therapeutic device during inhalation is that the cross-sectional area of the hose section is often narrowed to such an extent that the air resistance is considerable or the hose section is completely closed, with the effect that no air can get into the patient's airways."}
{"text": "1. Field\nMethods and apparatuses consistent with exemplary embodiments relate to detecting an error in gesture recognition.\n2. Description of the Related Art\nA dynamic vision sensor (DVS) may detect only an outline of a moving object, because the DVS responds to a change in light. Accordingly, high-speed processing is possible due to a relatively small quantity of data, and privacy may be maintained due to detection of the outline of the object. Also, low power may be consumed. However, there is a disadvantage in that, when only the outline is detected, it is difficult to accurately track a body part of a user, which may lead to an error in recognition.\nWhen a location of a user or a location of a hand for gesture recognition is not exactly found, there is another disadvantage in that a corresponding graphical user interface (GUI) may not properly operate, which may be confusing to a user who expects a control through gesture recognition."}
{"text": "From U.S. Pat. No. 7,127,049, a system enhancing automation of activating network service between a customer modem and a central office modem over a digital subscriber line link is known. In the system according to this prior art, the central office modem couples the customer modem to a network for providing the network service, the system comprising a polling system coupled with the central office modems.\nFurthermore, TS 33.203 of the third Generation Partnership Project (3 GPP), Access security for IP-based services (Release 10), 16 Jun. 2010 discloses a method for establishing a connection between an access node of a telecommunications network and a User Equipment.\nSuch known systems have a number of drawbacks. For example, due to the polling system defining a certain time period during which no network access is possible for the customer modem. Furthermore, the known system relays on the provisioning of a session assigned Internet Protocol address which cannot be used permanently and therefore necessitates the reconnection of the connection between the customer modem and the communications network (by possibly another Internet Protocol address and thus possibly a disconnection step and/or the initiation of new session and/or a reboot operation of the customers equipment.\nFurthermore, according to the prior art, in order to establish the Internet Protocol connection between, on the one hand, the Internet Protocol Edge node and, on the other hand, the User Equipment like a customer modem or a CPE (Customer Premises Equipment), it is always necessary to use—as an authentication information—an information that is used or distributed in an untrusted environment. For example, CPE units are pre-configured in a user-specific manner and distributed to a plurality of customers or the customer hast to configure the CPE with credentials he got from the operator of the telecommunications network. Such pre-configuration information is not inherently secure or trusted because, due to the steps occurring prior to the establishment of an Internet Protocol session according to the prior art, the telecommunications network operator necessarily needs to distribute such credentials or pre-configuration or other (previously trusted) information in an untrusted environment or to an untrusted environment (e.g. customers household).\nThese limitations have the effect that the connection between the customer modem on the one hand and the telecommunications network on the other hand is comparably time consuming. The user has to type in credentials to access the network. There is no plug and play solution to access the network and services provided to the network connectivity.\nFurthermore, German patent publication DE 10 2007 039 516 A1 discloses a method for configuring a communication port in a user-specific manner, the method comprising the step of providing a default profile, the default profile being directed to a specific user, and the method further comprising the step of assigning the user-specific default profile to a user-specific configuration profile, the user-specific configuration profile being assigned to a specific user."}
{"text": "Camera modules can be made compact and energy efficient for use not only in digital cameras but also in other devices such as mobile telephones, personal digital assistants (PDAs), and laptop computers. In a mobile telephone, for example, the camera module performs imaging functions, and components such as a keypad, a controller, and a display are used for operation of both telephone functions and imaging functions of the camera module. In particular, the display that provides a visual user interface for telephone configuration or text messaging can show an image from the camera module, and the keypad and controller that are used for making telephone calls can also control and activate the camera module. However, the camera modules in such applications have generally been limited to imaging functions.\nA common concern in portable devices, with or without camera modules, is power consumption. Components such as LCD displays or illuminated keypads, for example, consume power when the backlighting is activated. To conserve power and increase the time between required recharging of batteries, many portable systems such as mobile telephones turn off the backlighting after a short period of user inactivity. Additionally, the backlighting level may be under user control to permit use of lower intensity backlighting to conserve power when ambient lighting conditions permit. For example, backlighting for a keypad can be decreased or turned off when the ambient light level is high.\nAn ongoing goal for portable devices is to increase or improve the functionality of the devices while still providing an adequate operating time before batteries must be recharged."}
{"text": "Referring to FIG. 1, a known fan structure 10 mainly comprises a casing 11, a bearing 12, a shaft 13, a fan blade 14, a stator 15, a silicon steel sheet 16 and a magnet 17. The casing 11 has a shaft seat 11a, the bearing 12 has an outer ring 12a fixed on the shaft seat 11a and an inner ring 12b. The shaft 13 penetrates and is fixed at the inner ring 12b of the bearing 12 and the fan blade 14 is fixed at one end of the shaft 13. The stator 15 is disposed on the shaft seat 11a of the casing 11, the silicon steel sheet 16 and the magnet 17 are fixed on an inside wall 14a of the fan blade 14. When the fan structure 10 is actuated, the inner ring 12b of the bearing 12, the shaft 13, the fan blade 14, the silicon steel sheet 16 and the magnet 17 start to rotate simultaneously. The known fan structure 10 is no longer applied for miniaturized or thinned electronic products that have been made in recent years because which has a thick and heavy structure entirely."}
{"text": "A continuous combustion device usually has a primary combustion zone and a secondary combustion zone. Ideally, from a combustion or pollution aspect, or both, the primary combustion zone fuel/air ratio should be kept as close as possible to an optimum value which may be constant over the operating range of the combustion device. This does not normally happen. A gas turbine engine used as a propulsion unit on an aircraft, for example, will operate in varying operative conditions for different thrust settings. When an aircraft is on the ground, the thrust setting is relatively low to permit stopping or taxiing. When the aircraft initiates a take-off, the thrust is typically increased to its maximum setting until the aircraft reaches a cruising altitude and then is tapered back to an intermediate setting for a normal cruising flight. However, the fixed geometry of the conventional continuous combustion device provides a range of primary combustion zone fuel/air ratios which can go from over-rich to over-lean when the operative conditions vary.\nIt is well-known that the constituent emissions from a combustion device exhaust are formed by diverse processes depending on different, or even opposite, conditions, and therefore, problems are experienced when attempts are made to compensate for the variations in the operative conditions of the continuous combustion device. For example, the nitric oxide formation rate depends essentially on the temperature in the primary combustion zone and the availability of dissociated or free oxygen. A early or accelerated admission of cooling or dilution air to the primary zone can quench the reaction and restrict nitric oxide formation to low levels. This procedure may, however, increase hydrocarbons, smoke and carbon monoxide formation due to incomplete combustion.\nIn a conventional continuous combustion device used in a gas turbine engine at full load, carbon monoxide and hydrocarbons are practically non-existent, whereas nitric oxide emissions are at their peak. A continuous combustion device optimized for full load pollutant emissions would have a leaner than normal primary zone fuel/air ratio, and its yield in hydrocarbons and carbon monoxide would be higher, whereas nitric oxides would be considerably reduced, such a combustion device would not be practical for a normal application in a gas turbine engine where the fuel/air ratio is varied over a wide range, especially its stability would be poor and the emissions of hydrocarbons and carbon monoxide emissions would be very high when the engine is idling.\nTo maintain those objectionable or harmful exhaust emissions from a gas turbine engine combustor at an acceptable level, prior art combustion devices have provided means for varying the distribution of air flow within a combustor and means for providing automization, premixing and substantial vaporization to maintain the primary combustion zone fuel/air ratio within a narrow range when the operative conditions vary. One example of reducing harmful emissions in all modes of engine operations is described in U.S. Pat. No. 3,952,501, entitled GAS TURBINE CONTROL, naming John A. Saintsbury as inventor and issued Apr. 27, 1976. Saintsbury suggests a longitudinally adjustable baffle that is used to control the direction of air flow into the combustor to effect a substantially optimum proportionate distribution of combustion air throughout the combustor at all power levels. The fraction of primary zone airflow will be gradually reduced as the power is decreased, holding the fuel/air substantially to the predetermined optimum value. This procedure reduces the production of carbon monoxide and unburned hydrocarbons at low power because combustion takes place at a more favourable fuel/air ratio. The nitric oxide production is inherently low at reduced power because of the lower temperature of inlet air to the combustor. Moreover, more cooling air is diverted into the secondary zone, whereby the hot gases could be more efficiently cooled.\nThe nitric oxide produced in gas turbine engines is produced in the combustion process where the highest temperature in the cycle normally exists. Therefore, one way to limit the amount of nitric oxide produced is to limit the combustion temperature. Experience has shown that it is not enough to just limit the average temperature because when fuel is burned as drops of liquid or a diffusion gas flame, the combustion proceeds at near the stoichiometric value and the local temperature is very high, thus producing excessive nitric oxide. To produce the lowest possible nitric oxide, thoroughly premixing all of the fuel and combustion air in a mixing chamber separate from the combustion chamber itself is suggested in U.S. Pat. No. 5,477,671, entitled SINGLE STAGE PRE-MIXED CONSTANT FUEL/AIR RATIO COMBUSTOR and issued to Mowill on Dec. 26, 1995. Mowill describes in his patent, a compressed air valve and a fuel valve both under the control of a controller, to provide a preselected lean fuel/air ratio mixture for introduction to the combustion zone of an annular housing. Compressed air conduits are used to channel a portion of the total compressed air flow to a premixer and the remainder to a dilution zone of the combustor, and a fuel conduit is used to deliver all of the fuel to the premixer.\nAnother example is described in U.S. Pat. No. 3,905,192, entitled COMBUSTOR HAVING STAGED PRE-MIXING TUBES and issued to Pierce et al. on Sep. 16, 1975. Pierce et al. describe, in this patent, a gas turbine engine having an annular combustor with a plurality of staged premixing tubes extending from the forward end thereof. Each tube directs a flow to the combustor through two concentric flow passages. A moveable tube section is arranged to direct all of the air through both flow passages or just through one passage. Fuel is directed into the staged premixing tube for mixing with air generally flowing through the central flow passage. Swirler vanes are provided in each of the flow passages to provide for rotation of air passing therethrough. The air flow proportion through the two concentric flow passages can be varied by the moveable tube section and, therefore, the fuel/air premixing ratio is adjusted.\nHowever, since the proportion of air entering through the outer flow passage into the primary zone decreases as the proportion of air entering through the central flow passage into the primary zone in a premixed condition increases, the total amount of air reaching the primary zone through the both flow passages cannot be significantly regulated and, in fact, finally affects the improvement of the combustion conditions in the primary combustion zone."}
{"text": "1. Field of the Invention\nThe present invention relates to the conversion of thermal energy to electrical energy, and electrical energy to refrigeration, and more particularly to a solid state thermionic converter using semiconductor diode implementation.\n2. Relevant Technology\nThermionic energy conversion is a method of converting heat energy directly into electrical energy by thermionic emission. In this process, electrons are thermionically emitted from the surface of a metal by heating the metal and imparting sufficient energy to a portion of the electrons to overcome retarding forces at the surface of the metal in order to escape. Unlike most other conventional methods of generating electrical energy, thermionic conversion does not require either an intermediate form of energy or a working fluid, other than electrical charges, in order to change heat into electricity.\nIn its most elementary form, a conventional thermionic energy converter includes one electrode connected to a heat source, a second electrode connected to a heat sink and separated from the first electrode by an intervening space, leads connecting the electrodes to an electrical load, and an enclosure. The space in the enclosure is either highly evacuated or filled with a suitable rarefied vapor, such as cesium.\nThe essential process in a conventional thermionic converter is as follows. The heat source supplies heat at a sufficiently high temperature to one electrode, the emitter, from which electrons are thermionically evaporated into the evacuated or rarefied vapor-filled interelectrode space. The electrons move through this space toward the other electrode, the collector, which is kept at a low temperature near that of the heat sink. There the electrons condense and return to the hot electrode via external electrical leads and an electrical load connected between the emitter and the collector. The flow of electrons through the electrical load is sustained by the temperature difference between the electrodes. Thus, electrical work is delivered to the load.\nThermionic energy conversion is based on the concept that a low electron work function cathode in contact with a heat source will emit electrons. These electrons are absorbed by a cold, high work function anode, and they can flow back to the cathode through an external load where they perform useful work. Practical thermionic generators are limited by the work function of available metals or other materials that are used for the cathodes. Another important limitation is the space charge effect. The presence of charged electrons in the space between the cathode and anode will create an extra potential barrier which reduces the thermionic current. These limitations detrimentally affect the maximum current density, and thus present a major problem in developing large-scale thermionic converters.\nConventional thermionic converters are typically classified as vacuum converters or gas-filled converters. Vacuum converters have an evacuated medium between the electrodes. These converters have limited practical applications.\nEmbodiments in a first class of gas-filled converters are provided with a vaporized substance in the interelectrode space that generates positive ions. This vaporized substance is commonly a vaporized alkali metal such as cesium, potassium and rubidium. Because of the presence of these positive ions, liberated electrons can more easily travel from the emitter to the collector. The emitter temperature in this type of conventional device is in part determined by the vaporization temperature of the substance that generates the positive ions. Generally, the emitter temperature should be at least 3.5 times the temperature of the reservoir of the positive ion generating substance if efficient production of ions is to be achieved in these conventional devices.\nEmbodiments in a second class of gas-filled converters are provided with a third electrode to generate ions, and the gas in the interelectrode space in these conventional devices is an inert gas such as neon, argon, or xenon. Although these converters can operate at lower temperatures, such as about 1500 K, they are more complex.\nTypical conventional thermionic emitters are operated at temperatures ranging from about 1400 K to about 2200 K, and thermionic collectors are operated at temperatures ranging from about 500 K to about 1200 K. Under optimum conditions of operation, overall efficiencies of energy conversion range from 5% to 40%, electrical power densities are of the order of 1 to 100 watts/cm2, and current densities are of the order of 5 to 100 A/cm2. In general, the higher the emitter temperature, the higher the efficiency, the power and current densities with designs accounting for radiation losses. The voltage at which the power is delivered from one unit of a typical converter is 0.3 to 1.2 volts, i.e., about the same as that of an ordinary electrolytic cell. Thermionic systems with a high power rating frequently include many thermionic converter units connected electrically in series. Each thermionic converter unit is typically rated at 10 to 500 watts.\nThe high-temperature attributes of thermionic converters are advantageous for certain applications, but they are restrictive for others because the required emitter temperatures are generally beyond the practical capability of many conventional heat sources. In contrast, typical thermoelectric converters are operable at heat source temperatures ranging from about 500 K to about 1500 K. However, even under optimum conditions, overall efficiencies of thermoelectric energy converters only range from 3% to 10%, electrical power densities are normally less than a few watts/cm2, and current densities are of the order of 1 to 100 A/cm2.\nFrom a physics standpoint, thermoelectric devices are similar to thermionic devices. In both cases a temperature gradient is placed upon a metal or semiconductor, and both cases are based upon the concept that electron motion is electricity. However, the electron motion also carries energy. A forced current transports energy for both thermionic and thermoelectric devices. The main difference between thermoelectric and thermionic devices is in the transport mechanism; ballistic and diffusive transport for thermionics and ohmic transport for thermoelectrics. Ohmic flow is macroscopically diffusive, but not microscopically diffusive. The distinguishing feature is whether excess carriers are present. In thermoelectrics, the carriers normally present are responsible for current. In thermionics, the current is due to putting excess carriers in the gap. A thermionic device has a relatively high efficiency if the electrons ballistically go over and across the gap. For a thermionic device all of the kinetic energy is carried from one electrode to the other. The motion of electrons in a thermoelectric device is quasi-equilibrium and ohmic, and can be described in terms of a Seebeck coefficient, which is an equilibrium parameter.\nIn structures with narrow barriers, the electrons will not travel far enough to suffer collisions as they cross the barrier. Under these circumstances, the ballistic version of thermionic emission theory is a more accurate representation of the current transport.\nSolutions to the foregoing problems have been sought according to the present state of the art by using vacuum converters or gas-filled converters. Attempts to reduce space-charge effects with vacuum converters have involved the reduction of the interelectrode separation to the order of micrometers. Attempts to reduce the same effects with gas-filled converters have led to the introduction of positive ions into the cloud of electrons in front of the emitter. Nevertheless, these conventional devices still present shortcomings such as those related to limited maximum current densities and temperature regimes.\nConsequently, there remains a need to provide a more satisfactory solution to converting thermal energy to electrical energy at lower temperature regimes with high efficiencies and high power densities."}
{"text": "1. Field of the Invention\nThe present invention relates to protection for semiconductor integrated circuits from over-voltage and over-current conditions caused by electrostatic discharge (ESD). The present invention relates more specifically to a single structure ESD protection circuit using active devices integrated therein.\n2. Description of Prior Art\nESD events, which are inevitable in integrated circuit (IC) manufacturing, assembly, and application, generate huge over-current and over-voltage pulses that can fatally damage electronic components. ESD failure emerges as the major reliability concern to the IC industry, accounting for up to thirty percent of field failures. Therefore, industrial standards generally require on-chip ESD protection circuits for all IC chips.\nIn general principle, an ESD protection circuit is off, or open in normal operation and is triggered on, or closed, during an ESD event to form a low impedance path to shunt the current surge and to clamp the pin voltage to a sufficient level, thereby protecting an IC part from being damaged. To form a complete ESD protection scheme, shown in FIG. 1, protection against ESD pulses of two opposite polarities, or pulse modes, i.e., positive and negative, along a path from each input/output (I/O) pin to both power supply, or VDD, and ground, or GND, as well as the path between VDD to GND are required to achieve all direction protection.\nVarious on-chip ESD protection designs are known. Traditional solutions generally offer low impedance discharging paths in one direction only while relying on a passive current path with much higher impedance for protection in the opposite direction. This means four circuits may be needed for each I/O pin and an additional circuit between supply and ground, if a high level of ESD protection is desired. This solution consumes significant chip area, e.g., up to thirty percent of the whole chip.\nOther proposed solutions, such as U.S. Pat. No. 5,739,998 to Wada, may use various combinations of parasitic diodes and SCRs with transistor areas to provide protection for the various paths and pulse modes. Parasitic resistance and capacitance may negatively affect IC performance, as well as offer reduced protection performance because of higher impedance and slower reaction times which may not be adequate for very high speed chips such as RF IC\"\"s. In short, there is great need for a single structure circuit covering all ESD paths and pulse modes utilizing only active devices while utilizing minimal chip area.\nIn Proceedings of IEEE 7th International Symposium on the Physical and Failure Analysis of IC\"\"s, pp. 151-155, 1999, an ESD protection scheme is disclosed offering improved performance in a compact area providing low impedance active discharge paths in both directions to shunt ESD pulses in two opposite polarities. It therefore requires only two circuits for each I/O pin and a separate ESD circuit for VDD to GND for full ESD protection as seen in FIG. 2. This scheme represents a further size reduction in silicon area consumption over previously known schemes.\nAn ESD protection structure according to the present invention forms a single circuit which has a N-well, or N-substrate, having three separate P diffusions (hereinafter xe2x80x9ctubsxe2x80x9d) with the first and third P-tubs having N+ and P+ diffusions for ohmic contact therein. A second P-tub interposed between the first and third tub has two N+ diffusions and one P+ diffusion therein. The I/O pin, or anode, is connected across each of the two N+ diffusions and the P+ diffusion of the second P-tub. The VDD, or supply, pin, cathode one, is connected across the N+ and P+ diffusions of the third tub. The ground pin, cathode two, is connected across the P+ and N+ diffusions of the first P-tub. The structure thereby forms six bipolar transistors and eight parasitic resistors and thereby forming a circuit of five thyristors for control of the ESD over-current/over-voltage conditions. The structure set forth herein preferably forms symmetrical upper and lower protection subcircuits. The upper and lower subcircuit have one PNP transistor and two NPN transistors and four resistors creating two thyristors in each protection subcircuit, or structure half, and one thyristor formed between the subcircuits. One protection subcircuit provides positive and negative pulse protection between the I/O pin and the VDD pin. The other protection subcircuit provides positive and negative pulse protection between the I/O pin and the ground pin. The shared thyristor provides both positive and negative pulse protection between the ground pin and the VDD pin. By being built on a single substrate and providing protection against both polarities of surge along all paths, the present invention provides maximum protection with minimal use of IC surface area. The present structure operates symmetrically and responds in the sub-nanosecond range, providing low holding voltage, low discharge impedance adjustable trigger voltage and high failure threshold voltages."}
{"text": "Many electronic apparatus typically include coil components. Especially for mobile devices, coil components may have a chip form and may be surface-mounted on a circuit substrate included in the mobile devices. As an example of the prior art, Japanese Patent Application Publication No. 2006-324489 disclosed a chip coil including a helical conductor that is embedded in a hardened insulating resin and at least whose one end is coupled to an external electrode. The helical direction of the conductor is arranged in parallel with the surface of a substrate on which the coil is mounted.\nAs another example, Japanese Patent Application Publication No. 2014-232815 disclosed a coil component including a resin insulator, a coil-shaped inner conductor provided inside the insulator, and an external electrode electrically coupled to the internal conductor. The insulator is made in a cuboid shape with the length L, the width W, and the height H, where L>W≥H. The external electrode includes a conductor provided at each end of a plane perpendicular to the height H direction of the insulator as viewed in the length L direction. The internal conductor has a coil axis that is parallel with the width W direction of the insulator.\nIn the above-mentioned prior arts, insulators and conductors are alternately layered in the height direction using a photolithography and/or plating technique in order to obtain the coil component.\nIn recent years, miniaturization of components advances and so too with conductors and their sectional areas included in the components. As to electronic components in which insulator is made of resin, effects of stress applied from the resin to conductor cannot be ignored. For example, when an insulator and a conductor are formed of layers, a contraction stress caused by a hardening process of the insulator may act on a bonded portion between an upper conductive layer and a lower conductive layer, which may result in damage of the conductor. For this reason, miniaturization of conductors has been difficult."}
{"text": "Heretofore, it has been common practice to use small-diameter work rolls to roll hard materials, such as a magnetic steel strip, a stainless steel strip, and a high-tension steel strip, with a high strip thickness gauge accuracy. Rolling mills using such small-diameter work rolls are configured such that horizontally split housings, namely, a top inner housing and a bottom inner housing are used to respectively support a top roll group, which supports a top work roll and includes rolls arranged in a clustered form, and a bottom roll group, which supports a bottom work roll and includes rolls arranged in a clustered form. Further, a drive-side outer housing and a work-side outer housing are used to support the top inner housing and the bottom inner housing.\nA cluster-type multistage rolling mill of such type is disclosed in Patent Literature 1, for example."}
{"text": "Conventionally, there is known a program recording device for recording TV broadcast programs. In the program recording device, setting in advance the broadcast date and time, and the channel of a program which the user would like to record allows for the user to reserve recording of the program. Further, there is also known a program recording device configured such that setting a certain keyword in advance allows for the user to reserve recording of a program associated with the keyword, without the need of setting the broadcast date and time, and the channel individually.\nThere is known an electronic program guide (EPG), which is one of the tools for TV viewing, as a technique of supporting program recording. In the electronic program guide, a list of TV broadcast programs is displayed on a screen. This makes it easy for the viewer to know when and what programs are broadcasted. Utilizing the electronic program guide makes it possible to know the information relating to programs.\nIt is often the case that TV programs introduce a variety of shops, facilities, travel information, sightseeing information, or local information. As a result, the places, restaurants, or hotels introduced on TV programs are often set as destinations in a car navigation device. In this case, the user is required to record on paper or memorize the places introduced on TV programs, and set the recorded or memorized places as destinations in a car navigation device when the user would like to go to the places introduced on TV programs. However, the user may forget to record the places introduced on TV programs on paper, or may take time in setting the places as destinations in a car navigation device.\nFor instance, patent literature 1 and patent literature 2 disclose a method, in which a keyword relating to specific position information is extracted from program information received as an electronic program guide at the time of route search in a car navigation device, and the keyword is set as a destination in the car navigation device.\nFurther, it is often the case that the user is interested in the shops, facilities, or sightseeing areas where the user has visited as destinations. Therefore, when these places are broadcast on TV programs, it is desirable to automatically record these places as TV programs in which the user is interested. For instance, patent literature 3 and patent literature 4 disclose a method, in which a route to a destination is searched, and a program to be recorded is determined based on information relating to the destination or information relating to the route.\nHowever, in the conventional destination setting method utilizing an electronic program guide, it is necessary to search an intended destination from the entirety of an electronic program guide of an enormous information amount. Further, in the method for determining a recording program based on destination or route information, a program which is not necessary to be recorded may be recorded, because destination or route information of a car may be used by an unspecified family member who drives the car."}
{"text": "A patient in a healthcare facility is typically fitted with some type of device, such as a bracelet or other device, in order to positively identify that patient while they are receiving services. Existing remote patient management systems frequently utilize one of a plurality of identification technologies, for example, either barcode or radio frequency identification (RFID) tags, to accomplish patient identification and selection. Both the hardware and software deployed when implementing such a system are typically dedicated to the operation of a selected type of identification device using a corresponding protocol. In those situations where the management system software application does support both protocols, such support is achieved by providing alternative implementations for the respective technologies and protocols. In other words, the system user determines the technology to be deployed for patient selection prior to deployment of a patient management solution. For software vendors this requirement limits the available universe of potential healthcare provider customers to only those that have an implementation targeted at the identification technology that has already been chosen by the potential customer.\nIn those situations where a healthcare provider chooses to switch identification technology after an original software application is deployed, a software upgrade is mandatory, or if the patient management solution does not support the newly selected patient identification technology, a new patient management solution is selected or designed and implemented. From the perspective of an application software provider, the customer either targets a specific patient selection technology or incurs the cost of an alternative implementation. A system according to the principles of the present invention addresses these issues and their related problems."}
{"text": "In the hydrogenation of hydrocarbon fractions, catalysts are used which comprise a hydrogenation component usually on a suitable support such as, for example, silica, alumina, or mixtures thereof. The hydrogenation component is customarily supplied by a transition metal or metals, such as metals from Group VIB, Group VIII and some Group IB metals. Particularly suitable metals include copper, nickel, tungsten, cobalt, and molybdenum. The catalysts are conventionally prepared by impregnation of the support with the solution or dispersion of a decomposable salt of the metal hydrogenation component, followed by calcination in air to convert the salt to metal oxide. The hydrogenation component can be employed as the metal, metal oxide, metal sulfide, or mixtures thereof.\nIt is frequently desirable to convert the metal to a sulfide when the catalyst is used to process a sulfur-containing feed.\nIn general, the art teaches that sulfiding or conversion of the metal hydrogenation component to its sulfide form, can be accomplished by contacting the catalyst with hydrogen sulfide or organic sulfur compounds mixed with hydrocarbons at elevated temperatures up to about 1200.degree. F. However, loss in catalyst activity can result if sulfiding is carried out at a high temperature. Consequently, it is customary to employ relatively low temperatures, i.e., in the range of 750.degree. F. or lower, to avoid sintering of the metal components and in order to meet metallurgical limitations of the process equipment. For catalysts which have metals such as nickel incorporated by impregnation, U.S. Pat. No. 3,325,396 and U.S. Pat. No. 3,324,045 teach that sulfiding temperatures should not exceed 750.degree. F.\nU.S. Pat. No. 3,519,556 discloses an activation treatment for hydroconversion catalysts to increase hydrocracking activity which comprises sequentially calcining and sulfiding a metal and fluoride-containing hydrogel catalyst at temperatures of at least 930.degree. F. Fluoride-containing catalysts, however, have the drawback that fluoride leaches from the catalyst during use and forms hydrogen fluoride, which is extremely corrosive."}
{"text": "Automated banking machines are known in the prior art. Automated banking machines are commonly used to carry out transactions such as dispensing cash, checking account balances, paying bills and/or receiving deposits from users. Other types of automated banking machines may be used to purchase tickets, to issue coupons, to present checks, to print scrip and/or to carry out other functions either for a consumer or a service provider. For purposes of this description any device which is used for carrying out transactions involving transfers of value shall be referred to as an automated banking machine.\nAutomated banking machines often have the capability of accepting deposits from users. Such deposits may include items such as envelopes containing checks, credit slips, currency, coin or other items of value. Mechanisms have been developed for receiving such items from the user and transporting them into a secure compartment within the banking machine. Periodically a service provider may access the interior of the machine and remove the deposited items. The content and/or value of the deposited items are verified so that a credit may be properly applied to an account of the user or other entity on whose behalf the deposit has been made. Such depositories often include printing devices which are capable of printing identifying information on the deposited item. This identifying information enables the source of the item to be tracked and credit for the item correlated with the proper account after the item is removed from the machine.\nMany automated banking machines accept deposits from users in envelopes. Because the contents of the envelope are not verified at the time of deposit, the user's account generally is not credited for the deposit until the envelope is retrieved from the machine and the contents thereof verified. Often this must be done by persons who work for a financial institution. Delays in crediting a user's account may be experienced due to delays in removing deposits from machines, as well as the time it takes to review deposited items and enter appropriate credits. If the deposited items include instruments such as checks, further delays may be experienced. This is because after the instruments are removed from the machine they must be presented for payment to the appropriate institution. If the instrument is not honored or invalid the depositing customer's account cannot be credited for the deposit. Alternatively in situations where a credit has been made for a deposited instrument that is subsequently dishonored, the user's account must be charged the amount of the credit previously given. In addition the user commonly incurs a “bad check” fee due to the cost associated with the institution having to handle a dishonored deposit. All of these complications may result in delays and inconvenience to the user.\nAnother risk associated with conventional depositories in automated banking machines is that deposited items may be misappropriated. Because deposited checks and other instruments are not cancelled at the time of receipt by the automated banking machine, they may be stolen from the machine and cashed by unauthorized persons. Criminals may attempt to break into the machine to obtain the items that have been stored in the depository. Alternatively persons responsible for transporting items from the machine or persons responsible for verifying the items may misappropriate deposited instruments and currency. Alternatively the handling required for transporting and verifying the contents of deposits may result in deposited instruments being lost. Such circumstances can result in the user not receiving proper credit for deposited items.\nTo reduce many of the drawbacks associated with conventional depositories which receive deposits in the form of envelopes or other items, automated devices that can read and cancel deposited instruments have been developed. Examples of such devices are shown in U.S. Pat. Nos. 5,540,425 and 6,554,185 which are owned by the Assignee of the present invention. Such devices are capable of reading the coding on checks or other deposited items. For example bank checks include magnetic ink coding commonly referred to as “micr.” The micr coding on a check can be used to identify the institution upon which the check is drawn. The coding also identifies the account number of the user and the check number. This coding commonly appears in one or several areas on the instrument. Reading this coding in the automated banking machine enables the machine operator to determine the source of checks or other instruments that have been presented.\nImaging devices may also be used in processing instruments. Such imaging devices may be used to produce data corresponding to an image of the item that has been deposited. This image may be reviewed to determine the nature of the deposited item, and along with the information that can be obtained from the coding on the instrument allows processing of the credit to the user much more readily. Automated instrument processing systems also may provide the capability of printing an indication that the check or other instrument has been deposited and cancelled after it has been received. This reduces the risk that the instrument will subsequently be misappropriated and cashed by unauthorized persons.\nWhile automated deposit accepting and processing devices provide many advantages and benefits, existing devices may also have drawbacks. One drawback is that instruments must often be precisely aligned for purposes of reading micr coding or other indicia which is included on the instrument. This commonly requires special mechanisms to precisely position and align the instrument with the reading devices included in the device. Sometimes this requires aligning the check against a stop, which may crease or tear the check. A further drawback associated with some existing devices is that they are required to turn and reorient the deposited instrument. The mechanisms for doing this can be complex. Such complex mechanisms may encounter reliability problems due to the precise tolerances that must be maintained. Further difficulty is added by the fact that instruments that are received may be creased, torn or soiled. Handling such items may be difficult. Instruments becoming jammed in such mechanisms may result in costly repairs and downtime.\nA further drawback associated with some imaging systems in automated banking machines is that it is not practical to transmit an image of a deposited instrument for review and analysis at the time it is received. This is because the time and bandwidth necessary to capture and transmit an image of the deposited instrument may be longer than desirable. Extended transaction times may discourage the use of the machine. A further drawback is that even when images may be transmitted sufficiently quickly, the operator of the system is required to invest in the resources necessary to analyze the transmitted image and make a determination as to whether the deposited item should be accepted as valid or not. Such capabilities may include employees who must review the image and determine whether the item is genuine by comparison to data or other information such as examples of the customer's signature. Alternatively automated systems may be provided for analyzing the image of the instrument or the data printed or typed thereon. Providing such capabilities may be costly for the systems operator. Advances in photocopy technology also may make it difficult for operators of such systems to distinguish between genuine items and reproductions. As a result even with carefully operated and administered systems there is a risk that deposited items which are not genuine may be accepted.\nCertain standardized techniques have been developed for automated banking machine systems. The electronic message flows and formats commonly used for ATMs for example do not include the capability of transmitting a document image as part of the standard message which requests that a deposit transaction be authorized. As a result it has been difficult to achieve real time check verification and cashing in widely distributed systems. Further, in some systems it is difficult to readily correlate an image file with the particular transaction with which the image file is associated.\nThere may also be concerns in some systems with regard to the integrity of system data. For example there are risks that electronic image data representative of check images may be modified in order to carry out fraudulent activities. Also in some situations it may be desirable to mark on a check that has been presented with permanent markings to reduce the risk that the check can be improperly renegotiated. However, such markings may render it impossible to prove the character of the original check as it was received.\nA further drawback associated with some automated banking machine systems is that they cannot be used by individuals who do not have bank accounts. Generally automated banking machines require that deposited items be credited to a user's existing account with a financial institution. The user generally has to wait several days before the deposited item is verified and credited to the account. If the user does not have sufficient funds in the account to make a withdrawal, the user must generally wait for the verification process to be completed before the money may be withdrawn. This makes the use of automated banking machines generally unsuitable for persons who do not have bank accounts and/or cannot wait several days for deposited items to be verified and credited to their account.\nA further drawback associated with some existing automated banking machine systems is that some operators of such systems may wish to retain the capability to accept deposits in the form of items such as envelopes as well as checks and other instruments. Providing two separate depositories may add considerable cost and complexity to the machine. While mechanisms which can accept both single sheet-like instruments as well as envelopes have been developed, such mechanisms are often complex and unreliable. The capability of accepting both types of deposits is difficult to achieve because deposited instruments and envelopes may have varying thicknesses. The thickness of deposited envelopes may also be nonuniform. This is particularly true when such deposited envelopes may include items such as folded sheets or coin. Such combined depositories may also suffer from having lower security capabilities than mechanisms which are designed to accept only one type of deposit.\nThere is also often a desire to accept other types of documents in automated banking machines. Such documents may include for example utility bills or other items or instruments associated with value, or a particular account with which the customer may associate value or a particular payment. Such instruments may have thicknesses and properties which correspond to neither conventional checks or deposit envelopes. In addition the two-dimensional size of such items may also vary. This presents challenges for reliably handling such items. It may also be desirable in some circumstances to be able to image items and instruments which are associated with a customer. For example in some circumstances it may be desirable to receive a customer's driver's license, social security card, immigration card or other document to verify the identity of the user. In some cases it may be desirable to receive documents such as checks and currency bills and to process them using a common mechanism. Current depository mechanisms often do not have the capability of reliably handling or imaging such items.\nThus there exists a need for a deposit accepting apparatus and system for use in connection with automated banking machines that has the capability of handling and imaging more types of items, which may do so more reliably and which can be used in connection with more types of transactions and systems."}
{"text": "1. Field of the Invention\nThe invention relates to an apparatus for magnetically recording and/or playing back wide-band signals, such as video signals, on a magnetic tape, and more particularly to an apparatus having a rotatable head disc on which is disposed a number of electromagnetic heads which are spaced at equal radial distances from the axis of rotation of the disc and which are uniformly spaced around the circumference of the disc, and a tape guide having concave guide members for guiding the magnetic tape in a locally transversely curved condition past a part of the circular path of said electromagnetic heads in a direction extending transversely of the plane of rotation of the head disc. 2. Description of the Prior Art\nIn known apparatus of this type, to obtain a transversely curved shape of the magnetic tape, a partial vacuum is produced between two concave guide members in a vacuum chamber, so that the tape is drawn into the curvature of the guide members. Therefore, it is necessary to provide such apparatus with a vacuum pump, which is connected to the vacuum chamber by suitable means.\nAlthough the provision of a vacuum pump presents few problems in connection with large professional apparatus, which are used in broadcasting studios for example, the situation is different in the case of apparatus intended for domestic use. These latter apparatuses are mass-produced and are equipped with a different system for guiding the magnetic tape past the head disc. In such apparatus the magnetic tape is locally moved helically around a two-part cylindrical drum in which the head disc is rotatably mounted. A slot is provided between the two drum sections to allow the heads of the head disc, which pass through the slot, to write obliquely extending tracks on the magnetic tape. However, a drawback of this type of apparatus is that, for a number of reasons, the amount of information which can be stored on the magnetic tape is smaller than the amount that can be stored with an apparatus of the first type, in which the tracks are written transversely of the magnetic tape. Another drawback is that bringing the magnetic tape into contact with the drum, and maintaining it in contact, demands a intricate loop-shaped lace-up, as well as fairly intricate means for moving the tape from the disengaged position to the operating position and vice versa."}
{"text": "1. Field of the Invention\nThe present invention relates to adapters to be utilized in conjunction with aerosol spray cans particularly of the type providing a pistol grip spray gun attachment.\n2. Description of Prior Art\nU.S. Pat. Nos. 2,803,383; 2,820,578; 2,868,421; 2,877,934; 2,884,166; 3,112,849; 3,237,809; 3,304,797; and 3,506,159 relate to pistol grip type adapters for use with aerosol spray cans.\nAs far as is known, all of the spray can holders and adapters in the prior art have certain limitations in stability and ease of use. U.S. Pat. No. 2,803,383 disclosed a holder manufactured of spring steel designed to grip the aerosol can around the raised bead which surrounds the plunger element. That adapter is secured to the cannister by a spring loop of steel or wire which surrounds the uppermost bead of the cannister. That adapter does not provide a secure grip upon the body of the cannister and grips the can at its uppermost point thereby leaving the entire weight of the can below the gripping point. U.S. Pat. No. 2,820,578 disclosed a similar adapter that utilized a flange support member to engage the uppermost bead of the cannister. That adapter was secured to the cannister by a flange partially surrounding the circumference of the uppermost bead on the cannister thereby leaving the majority of the weight of the cannister below the gripping point of the adapter and failing to provide a secure grip upon the cannister.\nU.S. Pat. No. 2,868,421 disclosed a pistol grip type adapter wherein the uppermost bead of the cannister was received in a rectangular sleeve member which surrounded half the perimeter of the raised bead. That adapter provided only a partial support for the uppermost bead of the cannister.\nU.S. Pat. No. 2,877,934 disclosed a pistol grip type adapter which was secured to the cannister utilizing flexible finger members which engaged the uppermost bead of the cannister.\nU.S. Pat. No. 2,884,166 disclosed a pistol grip type adapter which utilized a clamping loop of spring steel to support the cannister from the uppermost bead. U.S. Pat. No. 3,112,849 disclosed a pistol grip type adapter similar to that of U.S. Pat. No. 2,877,934 embodying only a change in the method of joining the handle member to the cannister support member and utilized similar resilient finger flanges to engage the uppermost bead of the cannister.\nU.S. Pat. No. 3,237,809 disclosed a pistol grip adapter for use with specially designed aerosol containers and supported them in an inverted position. U.S. Pat. No. 3,304,797 disclosed a pistol grip type holder which was mounted to the cannister utilizing a resilient flange which snapped around the uppermost bead of the cannister. U.S. Pat. No. 3,506,159 disclosed a spray can holder which was mounted to the cannister by a resilient flange member which snapped onto the uppermost bead of the cannister.\nAs far as known, all of the prior art spray can adapters provide an inadequate grip upon the cannister so that the user cannot shake the cannister or move it quickly without fear of disengaging the adapter from the cannister. All of the known prior patents are designed to grip the cannister at the uppermost raised bead which immediately surrounds the plunger element in the cannister. Gripping the cannister at this point leaves the majority of weight and mass of the cannister below the gripping point, making the composite assembly of the cannister and the holder unbalanced and bottom heavy."}
{"text": "A CMOS imager circuit includes a focal plane array of pixel cells, each one of the cells including a photosensor, for example, a photogate, photoconductor or a photodiode overlying a substrate for accumulating photo-generated charge in the underlying portion of the substrate. Each pixel cell has a readout circuit that includes at least an output field effect transistor formed in the substrate and a charge storage region formed on the substrate connected to the gate of an output transistor. The charge storage region may be constructed as a floating diffusion region. Each pixel may include at least one electronic device such as a transistor for transferring charge from the photosensor to the storage region and one device, also typically a transistor, for resetting the storage region to a predetermined charge level prior to charge transference.\nIn a CMOS imager, the active elements of a pixel cell perform the necessary functions of: (1) photon to charge conversion; (2) accumulation of image charge; (3) resetting the storage region to a known state; (4) selection of a pixel for readout; and (5) output and amplification of a signal representing pixel charge. The charge at the storage region is typically converted to a pixel output voltage by the capacitance of the storage region and a source follower output transistor.\nCMOS imagers of the type discussed above are generally known as discussed, for example, in U.S. Pat. Nos. 6,140,630, 6,376,868, 6,310,366, 6,326,652, 6,204,524 and 6,333,205, assigned to Micron Technology, Inc., which are hereby incorporated by reference in their entirety.\nFIG. 1 illustrates a block diagram for a CMOS imager 10. The imager 10 includes a pixel array 20. The pixel array 20 comprises a plurality of pixels arranged in a predetermined number of columns and rows. The pixels of each row in array 20 are all turned on at the same time by a row select line and the pixels of each column are selectively output by a column select line. A plurality of row and column lines are provided for the entire array 20.\nThe row lines are selectively activated by the row driver 32 in response to row address decoder 30 and the column select lines are selectively activated by the column driver 36 in response to column address decoder 34. Thus, a row and column address is provided for each pixel. The CMOS imager 10 is operated by the control circuit 40, which controls address decoders 30, 34 for selecting the appropriate row and column lines for pixel readout, and row and column driver circuitry 32, 36, which apply driving voltage to the drive transistors of the selected row and column lines.\nEach column contains sampling capacitors and switches 38 associated with the column driver 36 reads a pixel reset signal Vrst and a pixel image signal Vsig for selected pixels. A differential signal (Vrst−Vsig) is produced by differential amplifier 40 for each pixel and is digitized by analog-to-digital converter 45 (ADC). The analog-to-digital converter 45 supplies the digitized pixel signals to an image processor 50, which forms a digital image output.\nLighting can effect image exposure. Light conditions may change spatially and over time. Thus, automatic light control is required to ensure that the best image is obtained by controlling the image sensor's exposure to the light. In some imager applications, there is a need to use the illumination during the actual exposure of an image (i.e., “present illumination”) to control the exposure (i.e., perform exposure control). That is, there is a need to use present illumination because the use of the previous picture's illumination may not be sufficient for the intended application.\nOne exemplary application that would benefit from using present illumination in exposure control is the imager in a swallowable pill application, such as the one described in U.S. application Ser. No. 10/143,578, now U.S. Pat. No. 7,187,049 issued Mar. 6, 2007, the disclosure of which is incorporated herein by reference. Due to the nature of the imager in a pill application, automatic light control using present illumination is required. A proposed solution would be to light the application's light source (e.g., light emitting diodes) prior to the actual exposure periods. This technique, however, creates an undesirable high waste of energy and power by having the light source on longer than the exposure period.\nAccordingly, there is a desire and need for automatic light control during an exposure period that uses present illumination, yet does not unnecessarily waste energy or power in the process."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical disk apparatus carrying out recording/reproduction of information with respect to an optical disk, and a method for controlling the same, and in particular, to an optical disk apparatus that automatically discriminates a standard of the optical disk.\n2. Description of the Related Art\nAs well known, in recent years, a technique of recording information at a high density has been developed, and an optical disk having a recording capacity of 4.7 GB (Giga Bytes) on one side has been put to practical use.\nAs this type of optical disk, there are various standards, for example, DVD-ROMs (Digital Versatile Disk—Read Only Memory) which are a reproduction-only type, DVD-RAMs (Random Access Memory), DVD-RWs (Rewritable), DVD+RWs which are a rewritable type, and DVD-Rs which a write-once type.\nIn an optical disk apparatus in which recording/reproducing can be carried out so as to selectively mount optical disks of several types of standards, there is provided a function in which the standard of a mounted optical disk is automatically discriminated by utilizing the fact that the light reflectance factors of the optical disks are different in accordance with the standard thereof, and adjustments or settings of respective portions are changed based on the discriminated result.\nIncidentally, the current rewritable type optical disks having the same standard are further subdivided into a plurality of types in accordance with a characteristic of a recording film to be used.\nTherefore, it is important for an optical disk apparatus to have a function of not only discriminating the basic standards of DVD-RAM, DVD-RW, DVD+RW, or the like, but also of automatically discriminating a characteristic of a recording film to be used, and of changing the adjustments and settings of the respective portions.\nIn Jpn. Pat. Appln. KOKAI Publication No. 2001-266367, there is disclosed a configuration in which focus-searching is applied to a mounted optical disk, and due to a light reflectance factor of the optical disk being measured based on a focus error signal obtained as a result thereof, a discrimination between a CD (Compact Disk)-ROM, a CD-R, and a CD-RW is carried out, and respective circuits for reproducing/recording are switched in accordance with a discrimination result.\nHowever, in the Jpn. Pat. Appln. KOKAI Publication No. 2001-266367, there is no description about that a type based on a characteristic of a recording film which is used is further automatically discriminated with respect to rewritable type optical disks having a same standard."}
{"text": "1. Field Of The Invention\nThis invention relates to computer systems and, more particularly, to methods and apparatus for controlling the output current produced by charge pumps used for erasing flash electrically-erasable programmable read only memory (flash EEPROM) arrays.\n2. History Of The Prior Art\nThere has been a recent trend toward lowering the power requirements of portable computers. In order to reduce power consumption, much of the integrated circuitry used in personal computers is being redesigned to run at lower voltage levels. The circuitry and components used in portable computers are being designed to operate at voltage levels such as five volts and 3.3 volts. This helps a great deal to reduce the power needs of such computers.\nUnfortunately, some features desired in portable computers require higher voltages. Recently, flash electrically-erasable programmable read o only memory (flash EEPROM memory) has been used to store basic input/output startup (BIOS) processes for personal computers. This flash EEPROM memory may be erased and reprogrammed without being removed from the computer by running a small update program when the BIOS processes are changed. However, erasing and reprogramming flash EEPROM memory requires approximately twelve volts to accomplish effectively, a voltage not available from the lower voltage batteries provided in personal computers.\nIn other electronic arrangements, charge pump circuits have been used to provide a high voltage from a lower voltage source. However, even though charge pumps have long been available which are capable of providing the voltages necessary for programming and erasing flash EEPROM memory arrays, no arrangement had been devised until recently for utilizing charge pumps integrated with flash EEPROM memory arrays to provide the voltages needed to accomplish erasing and programming of the flash EEPROM memory arrays using those positive source erase techniques which are used when twelve volts is available from an external source.\nThe primary reason for the failure is the universal perception that insufficient current can be generated using charge pumps to accomplish the erase process. The positive source method of erasing flash EEPROM memories draws a very substantial amount of current. However, recently it was discovered that using specially designed charge pumps sufficient current could be generated to accomplish positive source erase of flash EEPROM memory arrays. A charge pump arrangement for accomplishing positive source erase is disclosed in detail in U.S. patent application Ser. No. 08/119,719, now U.S. Pat. No. 5,414,669 entitled Method And Apparatus For Programming And Erasing Flash EEPROM Memory. Arrays Utilizing A. Charge Pump Circuit, K. Tedrow et al, filed on even date herewith, and assigned to the assignee of the present invention.\nOne of the problems encountered in providing charge pumps for generating source voltages for flash EEPROM memory arrays occurs because some manufactures desire to utilize the flash EEPROM arrays in circuitry for which an external power supply is available. This means that the memory array circuitry for providing the source voltage must be capable of utilizing power provided either by its internal charge pumps or by am external source. Typically the current available from an external power source is more than sufficient to erase the memory cells of the array during an erase operation. When provided by an external power source, the amount of current actually available is such that the memory cells are typically moved into a soft breakdown region in which current is dissipated through the source-substrate diode. The current transferred through the source-substrate diode junction is wasted.\nHowever, the current available from the charge pumps is less abundant. There is insufficient current available from charge pumps to erase the memory cells in a time period which is competitive with other forms of memory. For this reason, in order to assure that sufficient current to erase and program the array, it is necessary to assure that current is appropriately utilized so that the dissipation of large amounts of current through the source/substrate diode junction does not occur."}
{"text": "1. Technical Field\nThe present invention relates to measurement of frequencies, and more particularly, to a frequency measurement device capable of detecting a very small change in frequency.\n2. Related Art\nFrequency measurement methods are known such as a direct count method (for example, see JP-A-2001-119291) wherein the number of pulses passing within a prescribed gate period (gate time) is measured, a reciprocal method (for example, see JP-A-5-172861) wherein an accurate pulse period is measured and a frequency is calculated from the reciprocal of time, and a method (for example, see U.S. Pat. No. 7,230,458) wherein frequencies are calculated by obtaining a ΔΣ modulation signal. Although the direct count method can be realized with a relatively small-scale circuit, it is necessary to increase the gate time in order to increase frequency resolution (for example, the gate time needed for obtaining a resolution of 0.1 Hz is 10 seconds). Although the reciprocal method can solve this problem, the circuit for measuring pulse intervals accurately becomes larger in scale as compared to that of the direct count method. The method wherein frequencies are calculated by obtaining a ΔΣ modulation signal has a limitation in obtaining a dynamic range and has to trade off dynamic range against measurement accuracy. Although the measurement accuracy can be improved by obtaining a high-order ΔΣ modulation signal or using a PLL, there is a problem in that an additional analog circuit is needed, thus increasing the circuit size.\nThe present applicant has proposed a frequency measurement device based on a short gate time count method (also referred to as a short gate count method or a short gate method) as a new method substituting the above-described methods (for example, see JP-A-2009-250807). The short gate count method is configured to perform counting (sampling) of a measurement signal repeatedly without being interrupted at a prescribed short gate time to remove (filter out) high-frequency components from a series of obtained count values. In this way, the short gate count method improves both time resolution and frequency resolution as compared to the method of the related art. Since a frequency measurement device using this method is configured by a counter circuit and a small-scale computation circuit, it is possible to easily realize a multi-channel implementation while suppressing an increase in the circuit size. Moreover, as the gate time decreases, namely as the sampling frequency increases, it is possible to improve both dynamic range and resolution.\nHowever, when the short gate count method is used, pattern noise occurs under certain conditions. The pattern noise increases as an operating point parameter which is obtained by the ratio of a measurement frequency (the frequency of a measurement signal) to a sampling frequency (the frequency of a gate signal) approaches a simple advantageous value. Therefore, the present applicant has proposed a method of decreasing the level of pattern noise by selecting a measurement frequency and a sampling frequency based on the distribution characteristics of an operating point parameter and a noise level (see JP-A-2010-085286).\nHowever, in the method of the related art, it is difficult to reduce the level of pattern noise when the amount of change in the measurement frequency is large."}
{"text": "FIG. 1 shows an example of a prior art vibrator support structure. In FIG. 1, a vibrator 9 is made by adhering piezoelectric elements 8a and 8b to two side surfaces 1a and 1b, respectively, that are adjacent but separated by a ridgeline of vibrator body 1 which vibrator body I has the overall shape of a rectangular column. Two support filaments 3 and 4, which are nearly rectangular U-shaped when seen from the front, are attached to ridgeline 2 which is opposite the aforesaid ridgeline separating the piezoelectric elements 8a and 8b. The feet of these support filaments 3 and 4 are attached on substrates 5c and 5d, respectively, which substrates 5c and 5d are positioned so as to be mutually separated on damper material 6 disposed on a support stand 7.\nOne end of each lead line 10a and 10b is connected to the piezoelectric elements 8a and 8b, respectively, while the other end of each lead line 10a, 10b is connected to terminals 11a and 11b, respectively. These terminals 11a and 11b, as well as ground terminal 11c connected to support filament 4, are connected to an electric circuit (not shown) via the respective lead wires 12a, 12b and 12c. When a drive a.c. voltage is applied from an electric circuit (not shown) via these lead wires 12a, 12b, 12c, 10a and 10b, and terminals 11a and 11b, to the piezoelectric elements 8a and 8b, vibrator 9 undergoes flexural vibration along the directions of the X-axis depicted in the drawing.\nIn a support structure of this sort, the main purpose of arranging damper material 6 over support stand 7 is to shield out external vibrations in order to prevent external leakage of the vibration of vibrator 9 and to remove the effects imparted to vibrator 9 by external vibration. The reason why substrates 5c and 5d are mutually separated on a single damper material 6 is to remove any restraint on the flexural vibration of vibrator 9 by freeing the back and forth displacement of the mountings of the support filaments 3 and 4 relative to vibrator 9.\nHowever, this sort of vibrator support structure has problems in that the support posture of vibrator 9 may become rather unstable because substrates 5c and 5d are mutually separated on damper material 6. That is, when lead wires 12a, 12b and 12c are bent over during the packaging of the vibrator 9 arranged on support stand 7 together with the electric circuit connected to lead wires 12a, 12b and 12c inside a specified case, stresses such as tension and contraction occur in lead wires 12a, 12b and 12c. The subsequent transmittal of these stresses to substrate 5d causes deformations in damper material 6 that directly supports this substrate 5d, which damper material 6 comprises a soft material such as rubber or sponge. As a result, if the substrate 5d subsides somewhat after the stresses are removed, its left end side slants at an angle .THETA..sub.1, as shown for example by the side view in FIG. 2(a). The leg 4b of one support filament 4 attached to this substrate 5d compensates its posture of support for vibrator 9 by swaying somewhat toward the other support filament 3 from the vertical posture (as shown by the virtual line in the drawing) so that the stability of flexural vibration in vibrator 9 is lost. Also, when substrate 5d subsides at its left end side as shown in frontal view in FIG. 2(b), the single support filament 4 attached to substrate 5 distorts its posture relative to the vibrator axis thereby substantially decreasing the stability of flexural vibration.\nThis invention resulted from investigations undertaken to resolve the problems of such support structures. The object of this invention is to provide a vibrator support structure that does not impart an effect on the vibrator support posture arising from bending lead wires and, consequently, is capable of maintaining the constant stability of the flexural vibration of the vibrator."}
{"text": "This invention relates to improved rigid, dual radially locked subsea wellhead housings used in offshore drilling.\nEarly subsea wellhead system designs remained essentially unchanged for many years, typically employing a single landing or support shoulder between the 183/4\" high pressure wellhead housing (referred to below as \"the HP housing\") and the 30\" conductor housing (referred to below as \"the LP housing\"). This type of load-transfer mechanism operated relatively well for many years, until field failures occurred in the early 1980's while operating in high current situations. Subsequent studies revealed that external loadings, such as those generated by drilling risers or platform tiebacks, caused bending stresses in the casing extension below the HP housing, ultimately resulting in a fatigue failure of the weld between the casing and the bottom of the HP housing body. Numerous studies indicated that the magnitudes of the bending stresses in the extension below the HP housing were affected by several factors, including: (1) the type of interface between the HP and LP housings; (2) the radial clearances between the HP and LP housings; (3) the cement level in the annulus between the conductor pipe and casing extension below the HP housing; (4) the size of the wellhead extension; and (5) the external load.\nAs a result of the field failures and subsequent studies, the interface between the HP housing and the LP housing underwent several design changes. The end goal of most of the design efforts has been to eliminate fatigue failure in the casing immediately below the HP housing by increasing the rigidity between the HP and LP housings so that external loads imposed on the HP housing are more directly and efficiently transferred to the LP housing. Most designs implemented dual, axially spaced, sockets and a beveled landing shoulder at the interface between the HP and LP housings. However, prior \"dual socket\" designs included relatively high degrees of taper on the opposing interface elements, and require significant axial (or downward) force to keep the tapered elements in proper relative position so that the radial component of the tapers is maintained during load transfer. In other words, if vertical motion is present in the dual socket designs, the tapered sections of the landing shoulders can become misaligned, providing an increased radial separation or clearance between the HP and LP housings, resulting in decreased radial rigidity and greater potential for fatigue-related failure.\nThus, in order to maintain radial rigidity in such designs, it is also necessary to maintain axial rigidity. Therefore, significant attention is directed in such designs to providing means for maintaining axial rigidity to eliminate vertical movement between the HP and LP housings, often resulting in the use of complicated wedges or unidirectional latches. Further, as a consequence of meeting the goals of increasing both radial and axial rigidity, recent designs have resulted in complicated installation and locking procedures, often requiring the use of significant amounts of preload and expensive and cumbersome tools and camming devices. Moreover, once installed and used in accordance with such prior designs, the housings were difficult to refurbish and reuse.\nAccordingly, there is a need for significantly improved radial lock designs for subsea housings that are effective simple to install, maintain and refurbish."}
{"text": "In oilfield applications, pump assemblies are commonly invoked for pumping fluids at high pressures downward from the well surface downhole to a wellbore. Such oilfield operations frequently involve hydraulic fracturing. For hydraulic fracturing, abbreviated “fracking” herein for convenience, typically huge volumes of water and millions of pounds of abrasive-containing fluid such as sand and other primarily abrasive, solid crush-resistant particulate fracking materials Such primarily abrasive fracking materials are termed “proppant”; the combination of sand and concomitant proppant abrasive materials being termed “frac sand” herein. Frac sand and frac fluid, primarily comprising water, are pumped through the wellbore under extreme pressure and into targeted hydrocarbon reservoir regions proximal thereto, to create side “fractures” within the underlying hydrocarbon formations. As is well known in the art, in order to create such fractures, frac fluid containing abrasive proppant should preferably be pumped at extremely high pressures not only to facilitate fracture-creation, but also to sustain the propped-open subsurface structures. As will be appreciated by those skilled in the art, these propped-open structures afford additional pathways for hydrocarbon deposits to be urged to flow from underground formations to the well surface, thereby enhancing well production.\nPrior art apparatus and methodology for transporting frac fluid containing frac sand suffer from several long-standing disadvantages. For instance, after being loaded onto either special-purpose pneumatic trucks or railcars, frac sand has typically not been well sealed from environmental incursions during transfers. As a consequence of such environmental incursions, the integrity of this frac sand has been seriously undermined whereupon significant degradation attributable to cumulative effects of abrasion and friction, and exposure to moisture and rain occur during material transfer operations. Since it has been difficult, if not virtually impossible, for a special-purpose railcar to be brought sufficiently close to well sites, prevalent contemporary practice has been for several material transfers to be prerequisite for ultimate delivery of an adequate quantity of frac sand to the intended well site so that fracking operations may be initiated.\nMoreover, pneumatic trucks have frequently proven to be unavailable, thereby causing unloading of frac sand and the like from such special-purpose trucks or railcars to be likewise frequently delayed. It should be evident to those skilled in the art that a consequence of these multi-level delays is that railcars remain idle while substantial railroad demurrage fees accumulate. As, unfortunately, has become common in the art, under current fracking protocol, workers unavoidably are exposed to silica particles released during pneumatic transfer of frac sand. Another seemingly unavoidable disadvantage of the prior fracking art is that frac sand is frequently transferred from railcars to silos for storage. To achieve this sand storage capability at the transload facility, substantial capital costs of silos construction are incurred. According to the present invention, such sand storage silos and the like are rendered unnecessary because frac sand would preferably be stored in situ in the container framework hopper. Ergo, such containers can be stored and stacked on the ground independently of any associated railcars and pneumatic trucks—in the absence of silos.\nAlso well established in the fracking art is that once special-purpose trucks or railcars have transported frac sand to the prescribed well site at which fracking operations will occur, these trucks or railcars are returned empty to the point of origin to be reloaded and to effectuate another delivery of frac sand to the same or to another prescribed well site. Such a procedure obviously is not cost-effective and, on the contrary, is wasteful since substantial expense is incurred transporting these trucks or railcars that are devoid of any solid materials or liquid cargo.\nAccordingly, what is needed in the art is an apparatus and associated methodology for improving logistics of transporting frac sand proximal to well sites to inherently avoid such several long-standing prior art limitations and disadvantages, and thus enabling temporarily storing and transporting of frac sand to well sites preferably via conventional flatbed truck trailers and conventional flatbed railcars, and for contemporaneously temporarily storing and transporting frac fluids, i.e., spent water and associated hydrocarbon liquids from the well site back to the frac sand supplier and/or back to a water recycle facility.\nAs will become clear to those skilled in the art, this need is fulfilled by embodiments of the present invention which contemplate novel application of suitably configured container frameworks, including, but not limited to, standard International Organization for Standardization (“ISO”) shipping containers, constructed with internal structures for supporting and strengthening the internal walls and floor thereof, and a valve apparatus preferably configured to efficiently and securely achieve prerequisite frac sand transfer from such suitably configured container framework to well site fracking operations, regardless of the remoteness and limited accessibility of a diversity of well site locations.\nIt is an object and advantage of the present invention to provide a well-sealed temporary storage and transport container to inherently avoid environmental incursions thereinto which tends to undermine integrity of frac sand because of degradation attributable to cumulative effects of exposure to moisture and rain, and to abrasion and friction during material transfer.\nIt is an object of the present invention to avoid undue delay of transfer of frac sand from suppliers' facilities to well site locations due to unavailability of special purpose pneumatic trucks and railcars.\nIt is another object of the present invention to reduce hazardous exposure during fracking operations that threatens both workers and the environment.\nIt is yet another object and advantage of the present invention to provide a dual modality storage and transport container with simplified bladder valves enabling loading frac sand at a frac sand supplier's location, transferring and unloading such stored frac sand at the well site, and then enabling spent frac fluid comprising primarily flow-back water and associated small amounts of liquid hydrocarbons, to be loaded thereinto, and transported to a recycle water facility."}
{"text": "The specification relates to creating and organizing events in an activity stream. In particular, the specification relates to capturing, searching and collaborating on one or more events where the data is unstructured, the data including images, ink and audio recordings.\nThere is an overabundance of information in the world. People try and keep track of information and tasks but memory is fallible.\nPeople write down reminders on physical paper or save a text file to a computer. The text file can be accidently deleted and the paper can be thrown away, or the text file and paper can be forgotten. Another example is of a person taking a picture of a whiteboard to capture the discussion he had with a colleague. The picture can be hard to find if the user has taken many pictures on the same device.\nIn addition, collaboration between several people may be even more difficult. If a person wants to share a shopping list with her roommate they cannot both easily buy groceries without buying the same thing twice. Another example is of planning a trip. If two people are planning a trip together, it is difficult to collaborate on all aspects of the planning, especially if the two people are in different locations."}
{"text": "Modern communication and computer networks provide users with ready access to rich sources of information. For example, many vendors provide online product information, discount coupons, purchasing information, etc. Accessing such information often requires no more effort than navigating to the vendor's web site, or otherwise using an application, web browser, or other “user agent” to communicate with a vendor's server.\nHowever, vendors and other information providers frequently require users to leave email addresses, postal addresses, telephone numbers, or other personally identifying information. Some companies use identifying information only to provide the information in which users have expressed an interest, and they do not make such information available to third parties. Not all companies exercise such restraint. For example, personally identifying contact information provided by a user to one company frequently ends up being shared or sold to other companies, one or more of which may use the information to spam the user, or to otherwise send unsolicited information of questionable interest to the user."}
{"text": "The invention relates to the field of semiconductor packaging and testing, and in particular to mounting an integrated circuit that has been tested on a board, wherein the adjacent connecting pins of the integrated circuit are in an inline arrangement.\nIn order to keep the dimensions small for the package of a chip with multiple connecting pins (e.g., 88 pins), the connecting pins are arranged in line as close as possible. However, a minimum spacing must still be maintained since the space between associated holes on a board for insertion of the connecting pins may not be arbitrarily reduced. As the spacing of the holes and conducting lines on the board becomes closer, the process of fabricating a board becomes more complex and expensive. In addition, there is the risk that close spacing of the holes on the board will cause the connecting pins to short out during the soldering process.\nOne technique for increasing hole spacing on the board is to configure the connecting pins of the package by bending some of the connecting pins or, for example, by inwardly offsetting every second connecting pin. However, the advantage of greater hole spacing on the board is counteracted by the disadvantages of requiring a special test socket to test the chips and requiring an expensive packaging mechanism. Specifically, the packaging mechanism must be designed such that the connecting pins are not bent out of their precise alignment.\nTherefore, there is a need for a method of testing and mounting the package on a board in such a way that simple and inexpensively produced test sockets are sufficient for testing, and similarly simple and inexpensively produced packaging mechanisms are sufficient for packaging, while at the same time achieving an offset arrangement for the connecting pins."}
{"text": "When a gearshift is carried out in an automated manual transmission, in which gearshift an interlocking shifting element is involved, and when the gearshift is carried out which involves an interlocking shifting element connected between two gear systems, tooth-on-tooth positions can occur at the interlocking shifting element which prevent the interlocking shifting element from closing. To carry out the gearshift such tooth-on-tooth positions have to be eliminated.\nTo eliminate such tooth-on-tooth positions at interlocking shifting elements, actuators are used. From the prior art various methods or procedures are known for eliminating a tooth-on-tooth position at an interlocking shifting element, these various methods or procedures differing in particular in the actuators involved in the tooth-on-tooth position.\nFor the elimination of a tooth-on-tooth position at an interlocking shifting element fundamentally conflicting criteria have to be fulfilled, namely on the one hand the elimination of the tooth-on-tooth position at the interlocking shifting element should take place as quickly as possible, and on the other hand it should take place as comfortably as possible, granted that an elimination of a tooth-on-tooth position at an interlocking shifting element is the more comfortable, the less the drive-train is affected by knocks.\nFrom DE 10 2006 046 605 A1 a method for eliminating a tooth-on-tooth position at an interlocking shifting element of a transmission is known, in which prior art the tooth-on-tooth position at an interlocking shifting element can be eliminated by using a driving clutch or transmission brake as the actuator.\nFrom EP 1 572 491 B1 another method for eliminating a tooth-on-tooth position at an interlocking shifting element of a transmission is known, in which prior art, to eliminate the tooth-on-tooth position a clutch is regulated selectively and during this it is checked whether the tooth-on-tooth position has actually been eliminated. If after an operating attempt a tooth-on-tooth position cannot be eliminated, then in a subsequent operating attempt a target parameter of the clutch is increased in order to eliminate the tooth-on-tooth position in the next operating attempt.\nFrom DE 10 2005 012 308 A1 a further method for eliminating a tooth-on-tooth position at an interlocking shifting element is known, such that in this method too an actuator is operated with a defined target parameter to eliminate the tooth-on-tooth position at the interlocking shifting element. In this case the target parameter is a clutch position, which is adapted as a function of a reaction signal from a transmission."}
{"text": "1. Field of the Invention\nThe present invention relates to wafer processing, and more particularly to improving the wafer processing using site-dependent (S-D) procedures and subsystems.\n2. Description of the Related Art\nCurrent manufacturing methodology and factory design used for integrated circuits require many tools located as stand-alone platforms or grouped in general areas, usually separated by 2000 feet or more. Facilities to run these tools must therefore also be widely distributed throughout the factory. Typical functions required by these platforms are substrate coating (Adhesion, BARC, TARC, Resist, Top Coat), bake (post apply bake and post exposure bake) imaging (exposure), metrology (overlay, critical dimension, defect and film thickness), pre and post exposure cleaning using in immersion processing, etch (defining the pattern in the underling thin films) and post etch clean-up (polymer and other byproduct removal). Technologies targeting sub 32 nm gate lengths sub will require many of these operations to be repeated to complete a single active layer of the semiconductor device i.e. double BARC, double or triple patterning, double or triple imaging, etc.\nThe required gate level defect density for 15 nm gate technology is going to be approximately 0.01/cm2 at 10 nm in size per ITRS 2005 roadmap. Critical dimension control will need to be about 0.6 nm (3sigma), post etch, for the gate element. No lithographic and etch process tool exists with these performance capabilities.\nThese advanced technologies will need real time, wafer-to-wafer upstream adjustment of the process to maintain acceptable device results. Defectivity requirements will demand less movement of wafers from tool to tool within the factory as those movements add defects and factory clean room cost.\nPlatforms in use today function as manufacturing “islands”. This does not afford the best CoO development or allow for optimum process control. No 300 mm track design today can meet 300 wafers per hour throughput as claimed possible by some exposure tool manufacturers."}
{"text": "The present invention relates to improvements in flexible containers, and in particular, to modifications of the zipper elements for an easy open container.\nRecloseable, flexible containers are well-known in the art. Such containers normally comprise a bag-like structure made from a folded web of thermoplastic film, and typically include interlocking male and female zipper elements fused, extruded or attached to the bag side walls. The bag-like structure is created when the thermoplastic film is folded, severed and sealed along its side edges. The zipper elements are positioned on the opposing bag side walls and close to seal the bag mouth.\nBecause recloseable, flexible containers are, by definition, accessible, positive means to completely seal the bag mouth have been developed to preserve freshness, prevent accidental opening, and inhibit intentional tampering. In turn, means to open such sealed containers have also been developed. For example, tear strips are known which comprise weakened lines for tearing across sealed webs to gain access to the mouth of the bag. Naito, U.S. Pat. No. 3,198,228 and Uramoto, U.S. Pat. No. 3,780,781 (the '781 patent) both show tear strips having such weakened regions.\nProblems have arisen, however, in the use of containers having weakened regions for tear lines. Such problems include consumers' inability to tear the thermoplastic material; the need, despite the weakened regions, to use scissors to open the container, obviating the value of the tear strip; and irregular tearing down the bag lips, which compromises the strength, shape and operation of the zipper elements. Modifications, such as the inclusion of perforations along the weakened lines of tearing, have enhanced the efficacy of tear strips, but eliminate the airtight and tamper-proof features desired in completely sealing the mouth in such containers.\nOther tear strips comprised of a thin thread or string of fiber placed within the container at its top are shown by Staller, U.S. Pat. No. 3,543,343, as a means to open sealed containers along a single tear line. The drawback of such tear strips is that the tear line may deviate from the top of the bag and propagate into the zipper elements, again compromising their strength, shape and operation. In an apparent effort to avoid such problems, Uramoto, in U.S. Pat. No. 3,827,427 (the '427 patent), defines two tear lines, one on each lip of the bag mouth, between pairs of ribs. Further development is needed, however, to improve the design of tear strips for the easy opening of sealed flexible containers of thermoplastic material."}
{"text": "The present invention relates to an improved semiconductor integrated circuit that can transmit data at high speeds with reduced power dissipation and also relates to a system including such a circuit.\nPresent-day semiconductor integrated circuits (like an IEEE 1394 physical layer LSI) are specially designed to transmit data from unit to unit at high speeds to implement a system that can process a tremendously large quantity of data as fast as possible to cope with multimedia applications, in particular. As a result of recent vigorous research and development, the state-of-the-art semiconductor integrated circuits can normally transmit data at as high a rate as 200 megabits per second (Mbps). Some integrated circuits have already doubled the transmission rate although those circuits are just samples up to now. Furthermore, an extremely-high-speed data transmission technique exceeding 1 gigabits per second (Gbps) has also been realized albeit experimentally.\nIn the IEEE 1394 physical LSI, data is sent as direct current from the transmitting end through a differential twisted pair of cables. A resistor is connected between the twisted pair, and a potential difference between the twisted pair, which is variable with the current flowing through the resistor, will be supplied to the receiving end. The IEEE 1394 physical layer LSI also adopts a data transmission technique termed \"DS-Link\". Specifically, in transmitting data, a pair of data lines and a pair of strobe lines are associated with each single port. The specifications of the IEEE 1394 standards are described in Draft Standard Vex. 8.4.\nAccording to the IEEE 1394 physical layer LSI technology, however, a considerable amount of current is always consumed at the output ports because the direct current is output through the twisted pair. Particularly when an integrated circuit is provided with an increased number of output ports, the increase in current dissipation poses a serious problem. Moreover, although data is supposed to change its level on either the data line or strobe line pair in accordance with the DS-Link technique, the direct current is always flowing through the pair with no data transitions, thus also increasing the amount of currents consumed."}
{"text": "In recent years the increase in consumption of silver halide photographic materials has led to a strong request for a reduction of the processing times.\nA reduction in the exposure dose and, at the same time a reduction of the processing time for silver halide medical imaging elements have been strongly demanded in medical circles.\nTabular silver halide grains are crystals possessing two major faces that are substantially parallel. The average diameter of said faces is at least three times the distance separating them (the thickness).\nSilver halide photographic emulsions containing a high proportion of tabular grains have advantages of good developability, improved covering power and increased useful adsorption of sensitizing dye per weight of silver due to the high surface area-to-volume ratio of the grains. The use of such emulsions in photographic elements is disclosed in U.S. Pat. Nos. 4,425,425, 4,425,426, 4,433,048, 4,435,499, 4,439,520, and other related patents.\nThe use of automatic processors for the rapid processing (i.e., for a processing of from 45 to 90 sec) of light-sensitive silver halide elements including tabular silver halide grains, in particular light-sensitive silver halide elements for photographic use, is known. Such elements generally include a support (usually provided with a very thin subbing layer) having coated on at least one side a silver halide gelatin emulsion layer coated in turn with a gelatin protective layer. These elements are transported through machine processing units (developing, fixing, washing and drying) by means of opposed or staggered rollers (as described, for example, in U.S. Pat. No. 3,025,779) which also have the function of squeezing liquid from the film prior to drying.\nIf rapid processing of a film takes place, several problems occur, such as an inadequate image density (i.e. insufficient sensitivity, contrast and maximum density), insufficient fixing, insufficient washing, and insufficient film drying. Insufficient fixing and washing of a film can cause reduced image quality and modification of the silver tone. In order to reduce the time taken by the element to pass through the processing machine to 0.5 to 2 minutes, as particularly required in rapid processing of photographic elements, the processing is performed at relatively higher temperatures, usually higher than 30.degree. C., preferably between 35.degree. C. and 45.degree. C., such as 38.degree. C., and the gelatin content of the silver halide emulsions is considerably reduced as compared to that of emulsions for manual processing.\nUnder such conditions, the physical and photographic properties of the elements processed in an automatic processor tend to be worse. With high temperatures and such low gelatin content, for instance, the intrinsic sensitivity to pressure of the silver halide grains gets higher and the elements processed in the automatic processor show repeated marks caused by the pressure of the transporting rollers. Such pressure marks look like higher density regions and reduce the image faithfulness.\nIn order to prevent pressure marking, various methods have been described in the art. To this purpose, U.S. Pat. No. 2,960,404 describes the use in the photographic elements of glycerine, ethylene glycol and the like, Japanese Pat. No. 5316/1972 describes the use of 1,4-cyclohexane dimethanol and the like, and Japanese Pat. No. 4939/1978 describes the use of trimethylol propane. Another possible method of preventing pressure marking is increasing the degree of hardening of the gelatin layers, in particular of the external protective layers. Also, as another method, photographic elements are known wherein an intermediate gelatin layer is interposed between the support and the emulsion layer. For example, U.S. Pat. No. 3,637,389 describes a rapid processing photographic element wherein gradation, density and sensitivity are improved by applying such an intermediate gelatin layer between the support and the emulsion layer.\nHowever, known methods of preventing pressure marking when used in photographic elements including tabular silver halide grains have proved less effective. In particular, when the hardening degree is increased to achieve a very low swelling index and to improve its resistance to pressure desensitization, the photographic characteristics are reduced. Accordingly, the problem still remains of preventing pressure marking in photographic elements including light-sensitive tabular silver halide emulsions.\nEP 238,271 discloses a silver halide photographic material comprising at least one hydrophilic colloidal layer on a support, showing a melting time of from 8 to 45 minutes, and a water content of from 10 to 20 g/m.sup.2 upon completion of the washing step. The material is preferably processed in a developing solution comprising indazole and benzotriazole derivatives. The preferred processing time is 45 sec.\nU.S. Pat. No. 4,647,528 discloses a method of increasing both covering power and scratching resistance by using a particular polymeric hardener in a photographic material comprising a support coated with at least one silver halide emulsion layer containing tabular silver halide grains with an aspect ratio higher than 5:1. This material is processed in developing and fixing solutions containing hardeners.\nU.S. Pat. No. 4,847,189 discloses a silver halide photographic material comprising a silver halide emulsion having tabular silver halide grains with an aspect ratio not lower than 5:1 and showing a melting time of from 8 to 45 minutes. Also disclosed is a method for processing this material in processing solutions containing hardeners in a preferred time of 45 seconds.\nEP 341,637 discloses a processing solution for the rapid processing of silver halide photographic material. The processing solution comprises a compound capable of reducing the molecular extinction coefficient at 624 nm of an aqueous solution of a specific carbocyanine sensitizing dye. The processing solution reduces the residual stain and improves photographic characteristics of the processed photographic material.\nEP 327,133 discloses a developing solution comprising a thiadiazole derivative. This compound prevent sludge and provides a processed photographic material having good photographic and physical qualities.\nEP 321,948 discloses a method to reduce the fixing time by employing a silver halide photographic material comprising a subbing layer which contains an aqueous polymeric latex.\nEP 428,455 discloses a developing solution substantially free of both aldehydic hardening agents and silver halide solvents and containing a pyrazolidinone derivative, for use in the development of fore-hardened radiographic films.\nHowever, when performing a super-rapid processing of less than 45 sec the above mentioned disadvantages can not be overcome by these techniques, and thus there is still the need for a method of processing a silver halide photographic material which shows good photographic and physical characteristics when processed in a super-rapid processing."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical image measuring apparatus for applying a light beam to an object to be measured, particularly a light scattering medium, and for measuring a surface form or inner form of the object to be measured by receiving a reflected light beam or a transmitted light beam to produce an image of a measured form. In particular, the present invention relates to an optical image measuring apparatus for measuring the surface form or inner form of the object to be measured by using an optical heterodyne detection method to produce the image of the measured form.\n2. Description of the Related Art\nIn recent years, attention has been given to an optical image measuring technique for producing an image of a surface or inner portion of an object to be measured using a laser light source or the like. This optical image measuring technique is not hazardous to human bodies in contrast to a conventional X-ray CT. Therefore, the development of applications in the medical field has been particularly expected.\nAn example of a typical method in the optical image measuring technique is a low coherent interference method (also called an optical coherence tomography or the like). This method uses the low coherence of a broad-band light source having a broad spectral width, such as a super luminescent diode (SLD). According to the method, reflection light from an object to be measured or light transmitted therethrough can be detected at superior distance resolution on the order of μm (for example, see Naohiro Tanno, “Kogaku” (Japanese Journal of Optics), Volume 28, No. 3, 116 (1999)).\nFIG. 8 shows a basic structure of a conventional optical image measuring apparatus based on a Michelson interferometer, as an example of an apparatus using the low coherent interference method. An optical image measuring apparatus 200 includes a broad-band light source 201, a mirror 202, a beam splitter 203, and a photo detector 204. An object to be measured 205 is made of a scattering medium. A light beam from the broad-band light source 201 is divided by the beamsplitter 203 into two, that is, reference light R propagating to the mirror 202 and signal light S propagating to the object to be measured 205. The reference light R is light reflected by the beam splitter 203. The signal light S is light transmitted through the beam splitter 203.\nHere, as shown in FIG. 8, a propagating direction of the signal light S is set as a z-axis direction and a plane orthogonal to the propagating direction of the signal light S is defined as an x-y plane. The mirror 202 is shiftable in a direction indicated by a double-headed arrow in FIG. 8 (z-scanning direction).\nThe reference light R is subjected to Doppler frequency shift through z-scanning when reflected by the mirror 202. On the other hand, the signal light S is reflected on a surface of the object to be measured 205 and inner layers thereof when the object to be measured 205 is irradiated with the signal light S. The object to be measured 205 is made of the scattering medium, so reflection light of the signal light S may be a diffusing wave having random phases including multiples cattering. The signal light propagating through the object to be measured 205 and the reference light that propagates through the mirror 202 to be subjected to the frequency shift are superimposed on each other by the beam splitter 203 to produce interference light.\nIn the image measurement using such a low coherent interference method, a difference in optical path length between the signal light S and the reference light R is within a coherence length (coherent distance) on the order of μm of the light source. In addition, only a component of the signal light S which has phase correlation to the reference light R interferes with the reference light R. That is, only a coherent signal light component of the signal light S selectively interferes with the reference light R. Based on such fundamentals, the position of the mirror 202 is shifted by the z-scanning to change the optical path length of the reference light R, so that a light reflection profile of the inner layers of the object to be measured 205 is measured. The object to be measured 105 is also scanned with the irradiated signal light S in an x-y plane direction. The interference light is detected by the photo detector 204 during such scanning in the z-direction and the x-y plane direction. An electrical signal (heterodyne signal) outputted as a detection result is analyzed to obtain a two-dimensional sectional image of the object to be measured 205 (see Naohiro Tanno, “Kogaku” (Japanese Journal of Optics), Volume 28, No. 3, 116 (1999)).\nAssume that an intensity of the reference light R and an intensity of the signal light S which are superimposed by the beam splitter 203 are given by Ir and Is, respectively, and a frequency difference between the reference light R and the signal light S and a phase difference therebetween are given by fif and Δθ, respectively. In this case, a heterodyne signal as expressed by the following expression is outputted from the photo detector (for example, see Yoshizawa and Seta “Optical Heterodyne Technology (revised edition)”, New Technology Communications (2003), p.2).\nExpression (1)i(t)œIr+Is+2√{square root over (IrIs)} cos(2πƒiƒt+Δθ)  (1)\nThe third term of the right side of the expression (1) indicates an alternating current electrical signal and the frequency f if thereof is equal to a frequency of beat caused from the interference between the reference light R and the signal light S. The frequency fif of an alternating current component of the heterodyne signal is called a beat frequency or the like. The first and second terms of the right side of the expression (1) indicate direct current components of the heterodyne signal and correspond to a signal intensity of background light of interference light.\nHowever, when the two-dimensional sectional image is obtained by the conventional low coherent interference method, it is necessary to scan the object to be measured 205 with a light beam and to successively detect reflection light waves from respective regions of the object to be measured 205 in a depth direction (z-direction) and a sectional direction (x-y plane direction). Therefore, the measurement of the object to be measured 205 requires a long time. In addition, it is hard to shorten a measurement time in view of measurement fundamentals.\nIn views of such problems, an optical image measuring apparatus for shortening a measurement time has been proposed. FIG. 9 shows a fundamental structure of an example of such an apparatus. As shown in FIG. 9, an optical image measuring apparatus 300 includes abroad-band light source 301, a mirror 302, a beam splitter 303, a two-dimensional photo sensor array 304 serving as a photo detector, and lenses 306 and 307. A light beam emitted from the light source 301 is converted into a parallel light flux by the lenses 306 and 307 and a beam diameter thereof is widened thereby. Then, the parallel light flux is divided into two, that is, the reference light R and the signal light S by the beam splitter 303. The reference light R is subjected to Doppler frequency shift through z-scanning with the mirror 302. On the other hand, the signal light S is incident on the object to be measured 305 over a broad area of the x-y plane because the beam diameter is widened. Therefore, the signal light S becomes reflection light including information related to the surface and inner portion of the object to be measured 305 in the incident area. The reference light R and the signal light S are superimposed on each other by the beam splitter 303 and detected by elements (photo sensors) arranged in parallel on the two-dimensional photo sensor array 304. Thus, it is possible to obtain a two-dimensional sectional image of the object to be measured 305 in real time without light beam scanning.\nAn apparatus described in K. P. Chan, M. Yamada, and H. Inaba, “Electronics Letters”, Vol. 30, 1753 (1994) has been known as such a non-scanning type optical image measuring apparatus. In the apparatus described in the same document, a plurality of juxtaposed heterodyne signals outputted from a two-dimensional photo sensor array are inputted to signal processing systems arranged in parallel to detect the amplitude and phase of each of the heterodyne signals.\nHowever, in such a constitution when the spatial resolution of an image is increased, it is necessary to increase a number of elements of the array. In addition, it is necessary to prepare a signal processing system including a number of channels corresponding to the number of elements. Therefore, it is supposedly hard to actually use the apparatus in fields that require a high resolution image, such as a medical field and an industrial field.\nThus, the inventors of the present invention have proposed the following non-scanning type optical image measuring apparatus in JP 2001-330558 A (claims and specification paragraphs [0044] and [0072] to [0077]). The optical image measuring apparatus according to this proposal includes a light source for emitting a light beam, an interference optical system, and a signal processing portion. In the interference optical system, the light beam emitted from the light source is divided into two, that is, signal light propagating through an examined object arrangement position in which an object to be examined is arranged and reference light propagating on an optical path different from an optical path passing through the examined object arrangement position. The signal light propagating through the examined object arrangement position and the reference light propagating on the different optical path are superimposed on each other to produce interference light. The optical interference system includes a frequency shifter, light cutoff devices, and photo sensors. The frequency shifter shifts a frequency of the signal light and a frequency of the reference light relative to each other. In order to receive the interference light in the optical interference system, the interference light is divided into two parts. The light cutoff devices periodically cut off the two divided parts of the interference light to generate two interference light pulse trains with a phase difference of 90 degrees therebetween. The photo sensors respectively receive the two interference light pulse trains. The photo sensors each have a plurality of detection elements which are spatially arranged and separately detect the interference light pulse trains. The signal processing portion combines the plurality of interference signals obtained by the photo sensors to generate signals of the signal light which correspond to respective points of interest of a surface or inner layers of the object to be examined which is arranged in the examined object arrangement position along the propagation path of the signal light.\nIn the optical image measuring apparatus, the interference light in which the reference light and the signal light interfere with each other is divided into two parts. The two parts of the interference light are received by the two photo sensors (two-dimensional photo sensor arrays) and respectively sampled by the light cutoff devices (shutters) disposed in front of both sensor arrays. A phase difference of π/2 is set between sampling periods of the two divided parts of the interference light. Therefore, an intensity of the signal light and an intensity of the reference light which compose background light of the interference light and phase quadrature components (sine component and cosine component) of the interference light are detected. In addition, an intensity of the background light included in outputs from both the sensor arrays is subtracted from the outputs of both the sensor arrays to calculate two phase quadrature components of the interference light. An amplitude of the interference light is obtained based on the calculation result.\nHowever, according to the optical image measuring apparatus described in JP 2001-330558 A (claims, specification paragraphs [0068] to [0084], and FIG. 1), the sampling frequency of the interference light is not automatically set, so the setting takes an effort. In addition, a sampling frequency corresponding to a beat frequency cannot be easily set in the optical image measuring apparatus described in Japanese Patent Application Number 2001-330558, so it is difficult to perform high-precision sampling with ease.\nAn optical image measuring apparatus for detecting interference light using three photo sensors has been proposed in, for example, Japanese Patent Application 2004-52195 by the inventors of the present invention. In view of a reduction in manufacturing cost and simplification of the structure, it may be desirable to perform the same measurement, particularly, alternating current component measurement using a small number of photo sensors.\nAn available image sensor such as a charge-coupled device (CCD) camera has been widely used for the two-dimensional photo sensor array of the optical image measuring apparatus as described above. However, up to now, a problem has been recognized that a currently available CCD camera cannot follow the beat frequency of a heterodyne signal which is on the order of several kHz to several MHz because of the low frequency response characteristic thereof. The feature of the optical image measuring apparatus which is proposed by the inventors of the present invention and described in JP 2001-330558 A is to perform the measurement using the low frequency response characteristic based on the sufficient understanding of the problem."}
{"text": "Certain types of prior personal computer systems each include a microprocessor (also referred to as a central processing unit or CPU) that is coupled to several types of storage systems--namely, a read only memory (\"ROM\"), a random access memory (\"RAM\") or dynamic random access memory (\"DRAM\"), a hard (i.e., fixed) disk drive for mass storage, and a floppy disk drive or drives for storage on removable magnetic floppy disks. The floppy disks are also referred to as diskettes.\nA relatively new mass storage device is a flash memory card. One prior art flash memory card includes flash electrically erasable programmable read only memories (\"flash EEPROMs\") and an electrical connector as part of a plastic package that is smaller than a 3.5 inch floppy disk. The flash memory card can be connected to a personal computer via the electrical connector.\nThe flash EEPROM is a nonvolatile memory that can be programmed by the user. Once programmed, the flash EEPROM retains its data until erased. Electrical erasure of the flash EEPROM erases the entire contents of the memory of the device in one relatively rapid operation. The flash EEPROM may then be reprogrammed.\nThat prior art flash memory card allows for the storage of data files and application programs on the purely solid-state medium of flash EEPROMs. System resident flash filing systems permit the prior art flash memory card to function as if it were a physical disk drive. The prior flash memory card in conjunction with a flash filing system provides an alternative to both a fixed hard disk and a floppy disk in a Disk Operating System (\"DOS\") compatible portable personal computer (\"PC\").\nThe storage of user application software on a prior art flash memory card substantially reduces the relatively slow prior art disk-DRAM download process. That prior art flash memory card can be read from faster than certain prior art hard disk drives. That prior art flash memory card generally uses less power than certain prior art hard disk drives. That prior art flash memory card is also smaller and lighter than certain prior art hard disk drives.\nPrior art personal computers typically have a redundant memory structure--i.e., there is a DRAM for storage of applications and data to be executed plus a hard disk or a floppy disk for mass storage. Applications and data need to be loaded into the DRAM.\nThe prior art flash memory card, however, has a read access time and a command register microprocessor write interface that permits an \"execute-in-place\" architecture. This configuration eliminates the need for a DRAM. Thus, redundancy associated with having both a DRAM and disk drive is eliminated.\nCertain prior art flash memory cards can be transported from personal computer to personal computer. Moreover, the flash EEPROMs of certain flash memory cards are nonvolatile and thus do not require a battery backup.\nOne type of flash EEPROM used in a prior art flash memory card has a standby mode that disables most of the flash EEPROM circuitry and reduces device power consumption. This flash EEPROM also has an active mode. The active mode requires increased power consumption. The active mode is used when the flash EEPROM is being written to, read from, or erased.\nOne disadvantage, however, of certain flash memory cards is that a microprocessor has to step through relatively complex erasure or programming software routines in order to erase or program the flash memory card.\nAnother disadvantage of certain flash memory cards is that even when the flash EEPROMs making up the card are in the standby mode, the flash memory card still consumes a significant amount of power. The amount of power consumed by a flash memory card with flash EEPROMs in the standby mode is especially noticeable when the flash memory card is used in conjunction with a battery-powered laptop personal computer. In order to extend battery life, power consumption must be kept to a minimum for a battery-powered laptop personal computer.\nRecently, improvements have been made in flash EEPROMs. Memory capacity has been increased. In addition, one type of prior art flash EEPROM includes a write state machine on the flash EEPROM chip.\nThe write state machine comprises circuitry that automatically steps the prior art flash EPROM through a multistep program or erasure sequence upon receiving an initiating command from a microprocessor. The prior art flash EPROM includes a ready/busy output pin that indicates whether the write state machine is ready to accept a command or whether the write state machine is currently busy programming or erasing the flash EEPROM.\nThat type of prior art flash EEPROM also includes a power down mode that can be initiated by applying a logical signal to a power down pin. In the power down mode, the flash EEPROM consumes less power than in the standby mode."}
{"text": "The present invention relates to a chemical mechanical polishing carrier head that includes a flexible membrane, and associated methods.\nIntegrated circuits are typically formed on substrates, particularly silicon wafers, by the sequential deposition of conductive, semiconductive or insulative layers. After each layer is deposited, it is etched to create circuitry features. As a series of layers are sequentially deposited and etched, the outer or uppermost surface of the substrate, i.e., the exposed surface of the substrate, becomes increasingly nonplanar. This nonplanar surface presents problems in the photolithographic steps of the integrated circuit fabrication process. Therefore, there is a need to periodically planarize the substrate surface.\nOne accepted method of planarization is chemical mechanical polishing (CMP). This planarization method typically requires that the substrate be mounted on a carrier or polishing head. The exposed surface of the substrate is placed against a moving polishing surface, such as a rotating polishing pad. The polishing pad may be either a “standard” polishing pad with a durable roughened surface or a “fixed-abrasive” polishing pad with abrasive particles held in a containment media. The carrier head provides a controllable load to the substrate to push it against the polishing pad. A polishing slurry, including at least one chemically reactive agent, and abrasive particles if a standard pad is used, is supplied to the surface of the polishing pad.\nSome carrier heads include a flexible membrane with a mounting surface that receives the substrate. A chamber behind the flexible membrane is pressurized to cause the membrane to expand outwardly and apply the load to the substrate. After polishing, the substrate is chucked to the mounting surface, lifted off the polishing pad, and moved to another location, such as a transfer station or another polishing pad."}
{"text": "The invention relates to a distributing device for cut tobacco, to be used in cigarette making machines, which includes a pointed rotating drum which takes or receives the cut tobacco, and a pointed detaching roller which cooperates with said drum and detaches the cut tobacco from same.\nThe pointed drum can take the cut tobacco from a bulky mass of tobacco contained in a box. In this case the pointed drum cooperates also with a pointed equalizer drum, rotating and preferably with a shorter diamater, and which removes the excess of tobacco taken by the drum, leaving, attached to the points of the latter, only a thin layer of cut tobacco, constant in thickness, and which is then removed from the drum by the action of the detaching cylinder.\nThe pointed drum can, however, receive the cut tobacco, already previously spread and carded which is fed from the top into a vertical duct, wherein a column of tobacco is formed. This tobacco is taken by the pointed drum, in the form of a uniform and thin veil, at the basis of the said vertical duct, and it is then removed from the drum by the detaching roller.\nIn the cut tobacco distributing devices of the known previous art, the rotating detaching roller is also provided with points and, therefore, it exerts a mechanical action on the tobacco, which damages the tobacco itself and determines a substantial degradation of same.\nThe invention has the purpose of eliminating this inconvenience and of ensuring a much more delicate treatment of the cut tobacco.\nThis purpose is achieved by the invention through the fact that the points of the drum are retractable inside the drum itself and are controlled by means which retract the same temporarily in a complete or almost complete manner in the region of the detaching cylinder, which is constructed as a smooth (without points) pneumatic suction cylinder. In this way the cut tobacco kept in the points of the drum is released by the complete and temporary retraction of such points in the region of the detaching cylinder which removes the tobacco from the drum by simple suction, without any sensible resistance and practically with no mechanical damage. Subsequently the points of the drum are caused to protrude again, to take or receive again a layer, or veil, of cut tobacco.\nAccording to one embodiment of the invention the means which control the travel of the retractable points in the drum are adjustable in such a way as to make it possible to vary the degree of projection of the points from the body of the drum, in order to vary the quantity of cut tobacco taken or received in the time unit, by the pointed drum itself."}
{"text": "1. Field of the Invention\nThe present invention relates to a combination lock, and more particularly to a combination lock having therein a knob rotatably mounted on the lock to activate/deactivate the locking mechanism of the lock.\n2. Description of the Prior Art\nIn public places, such as the train station, the Gray Hound station or the Supermarkets and the malls, a large number of storage boxes are installed for users to temporarily store therein personal items. This kind of storing boxes normally has a built-in unlocking mechanism so that when the key to the storing box is lost or for any unknown reasons the key to the specific storing box is lost, the supervisor may always unlock the storing box. However, there is no way that the responsible authority is able to monitor all the users of the storing boxes, so that some users may copy the keys to the storing boxes when and even after they rent the storing boxes. Then these people may use the copied keys to unlock the storing boxes and steal whatever the items received inside the storing boxes.\nTo overcome the shortcomings, the present invention tends to provide an improved lock to mitigate the aforementioned problems."}
{"text": "A standard apparatus for erecting box blanks comprises a supply holding a stack of flattened box blanks, a conveyor passing adjacent the supply and having a succession of box seats, and a planet carrier rotatable about a sun axis between the supply and the conveyor. Several planetary elements on the carrier offset from the sun axis have axially spaced inner and outer ends and are rotatable about a planet axis parallel to and offset from the sun axis. Grabs on the planetary-element outer ends are engageable with the blank in the supply and engageable in the conveyor seats. A planet drive gear fixed to the planetary-element inner end is rotatable with the planetary element about the planet axis.\nAccording to German 42 24 897 a differential drive including a sun wheel rotatable about the sun axis and meshing with all of the planet gears moves the planetary element and the carrier such that the grabs follow a hypocycloidal path having for each grab a respective outer point. Thus each grab engages the end flattened blank in the supply, pulls it out, then swings it through an arc to deposit it in one of the seats or cells of the conveyor and then swings through a further arc to return to the starting position. During this movement the blank is opened up or erected into its three-dimensional shape.\nThe problem with such a system is that the large symmetrical paths followed by the plural grabs take up a great deal of space. Several grabs are needed, however, to keep up with production pace, since each grab moves through a long complex path between its actual end-point use positions. In addition resetting the machine to work with box blanks of different format is a fairly complex task that takes quite some time, as each of the grabs has to be adjusted and reset. The grabs driven from a common sun gear have to be set perfectly so that each does its function when it moves into position between the supply and the conveyor.\nAccordingly U.S. Pat. No. 4,518,301 of Greenwell shows a system for taking flattened box blanks out of a supply and transferring them to a conveyor. The blanks are erected as they are moved between the supply and the conveyor by a system having several grabs whose angular positions are controlled by a complex cam system. Such an arrangement is also extremely difficult to reset when the size of the blanks being handled changes."}
{"text": "In modern internal combustion engines, camshaft adjusting systems are used for the variable adjustment of the control times of gas exchange valves, in order for it to be possible for the phase relation between the crankshaft and the camshaft to be designed in a defined angular range in a variable manner. Hydraulic camshaft adjusting systems usually require energy for the adjusting operation, which energy is made available to them via a pressure medium system. Here, the energy which is available in the pressure medium system is formed from the product of pressure medium pressure and pressure medium volumetric flow. During an adjustment of one or more camshafts of the internal combustion engine, the required pressure medium flow, which flows into the camshaft adjuster, is branched off by means of at least one control valve which is connected upstream of the camshaft adjuster. At the same time, an equally large pressure medium flow passes out of the camshaft adjuster back into the pressure medium reservoir of the pressure medium system, and is therefore fed to the pressure medium system again. As a result of the removal of the pressure medium flow out of the pressure medium system, the pressure medium pressure falls as a consequence. As a result, the energy is also reduced which is available to the camshaft adjuster for the adjustment. In order to prevent this, a pressure accumulator can be provided which makes pressure energy available to the camshaft adjusting system in the case of a drop in the pressure medium pressure, for equalization purposes.\nA volume accumulator which communicates with a camshaft adjuster is described in a patent application which has already been filed by the applicant under the reference number DE 10 2009 049 459.6 but has not yet been published. The volume accumulator is formed in a cavity of a camshaft and is connected, via a pressure medium channel which is arranged in the camshaft, to a control valve which serves to supply the camshaft adjuster with pressure medium. In this embodiment, the circumstance has a disadvantageous effect that the control valve and volume accumulator are configured separately. As a result, said components have to be mounted individually, in each case on one side of the camshaft, and are connected to one another by way of a pressure medium channel in the camshaft. In addition, throttling losses can occur in the pressure medium channel, in particular because of cross-sectional constrictions, additional baffles or deflection means, as a result of which the function of the pressure accumulator is impaired."}
{"text": "The present invention generally relates to perpendicular magnetic recording and reproducing systems, and more particularly to a perpendicular magnetic recording and reproducing system which realizes satisfactory perpendicular magnetic recording and reproducing characteristics.\nGenerally, when recording and reproducing a signal on and from a magnetic recording medium by use of a ring core head as a magnetic head, the ring core head magnetizes a magnetic layer of the magnetic recording medium in a longitudinal direction of the magnetic recording medium (that is, in an in-plane direction) at the time of the recording and picks up the recording at the time of the reproduction. However, according to such a longitudinal magnetic recording system, it is known that the demagnetization field becomes high as the recording density increases and the demagnetization field introduces undesirable effects to the high density magnetic recording. Hence, in order to eliminate the undesirable effects of the demagnetization, a perpendicular magnetic recording system has been proposed in which the ring core head magnetizes the magnetic layer of the magnetic recording medium in a direction perpendicular to the magnetic layer. According to the perpendicular magnetic recording system, the demagnetization field becomes low as the magnetic recording density increases, and theoretically, it is possible to realize a satisfactory high density magnetic recording in which there is no decrease in the remanent magnetization.\nAs a conventional perpendicular magnetic recording medium which is used in the perpendicular magnetic recording system, there is a perpendicular magnetic recording medium having a cobalt-chromium (Co--Cr) film formed on a base film by a sputtering process. As is well known, the Co--Cr film is extremely suited for use in the perpendicular magnetic recording medium because the Co--Cr film has a relatively high saturation magnetization (Ms) and favors magnetization in a direction perpendicular to the Co--Cr film (that is, the coercivity in the direction perpendicular to the Co--Cr film is large and the axis of easy magnetization is perpendicular to the Co--Cr film).\nHowever, when the ring core head performs the perpendicular magnetic recording and reproduction with respect to the perpendicular magnetic recording medium having the sputtered Co--Cr film, it is impossible to concentrate the magnetic flux at a predetermined magnetic recording position on the perpendicular magnetic recording medium, and there is a disadvantage in that it is impossible to obtain a strong magnetization which is in the direction perpendicular to the Co--Cr film and does not spread in the longitudinal direction of the perpendicular magnetic recording medium. In other words, when the ring core head is used to perform the recording on the Co--Cr film of the perpendicular magnetic recording medium, the magnetization direction easily deviates in the longitudinal direction of the perpendicular magnetic recording medium since the magnetic field generated by the ring core head includes considerable components in the in-plane direction. Accordingly, in order to maintain the magnetization direction in the perpendicular direction, the perpendicular magnetic recording medium must have a high perpendicular anisotropic magnetic field and have a saturation magnetization which is suppressed to a certain extent. However, the Co--Cr film does not have such characteristics, and there is a disadvantage in that it is impossible to perform a satisfactory perpendicular magnetic recording by the ring core head. In addition, the coercivity in the perpendicular direction must be large in order to obtain a high reproduced output from the perpendicular magnetic recording medium having the Co--Cr film. On the other hand, it is desirable to make the thickness of the perpendicular magnetic recording medium large in order to decrease the demagnetization field, but the perpendicular magnetic recording medium will not make contact with the ring core head in a satisfactory state when the thickness of the perpendicular magnetic recording medium is large because the perpendicular magnetic recording medium will lose its flexibility and become rigid. In this case, there are disadvantages in that the rigid perpendicular magnetic recording medium is easily damaged and undesirable effects are introduced to the ring core head, and it is impossible to perform a satisfactory perpendicular magnetic recording and reproduction.\nAccordingly, a perpendicular magnetic recording medium having a double film construction has been proposed. According to this perpendicular magnetic recording medium, a film having a high permeability, that is, a film having a low coercivity such as a nickel-ion (Ni--Fe) film, is formed between the Co--Cr film and the base film. The magnetic flux which is spread within the high permeability film is concentrated toward the magnetic pole of the ring core head at a predetermined magnetic recording position in order to obtain a strong magnetization which is in the perpendicular direction and does not spread in the longitudinal direction of the perpendicular magnetic recording medium. However, in the case of the perpendicular magnetic recording medium having the double film construction, the coercivity of the high permeability film is extremely small compared to the coercivity of the Co--Cr film, and there is a disadvantage in that Barkhausen noise is generated. For example, the coercivity of the Co--Cr film is over 700 Oe, and the coercivity of the high permeability film is under 10 Oe. Further, in order to produce the perpendicular magnetic recording medium having the double film construction, an amorphous (ion-nickel) Fe--Ni alloy or the like is formed on the base film by a sputtering process under a predetermined sputtering condition suited for forming the high permeability film, and Co--Cr is thereafter formed on the high permeability film by a sputtering process under a certain sputtering condition suited for forming the Co--Cr film. As a result, the sputtering condition under which the sputtering process is performed and the target must be changed for the formation of each film, and the sputtering processes cannot be performed continuously. Therefore, there are disadvantages in that the processes of manufacturing the perpendicular magnetic recording medium are complex and unsuited for mass production.\nOn the other hand, as the magnetic which is used in the perpendicular magnetic recording and reproducing system, there are in addition to the ring core head an auxiliary pole type (or driven) perpendicular magnetic head (hereinafter simply referred to as an auxiliary pole type head) having an auxiliary magnetic pole opposing a main magnetic pole and a single-sided main pole type (or driven) perpendicular magnetic head (hereinafter simply referred to as a single-sided main pole type head) not requiring an auxiliary magnetic pole and having a main magnetic pole opposing only one surface of the perpendicular magnetic recording medium. However, in the case of the perpendicular magnetic recording and reproducing system using the auxiliary pole type head, the perpendicular magnetic recording medium must be inserted between the main and auxiliary magnetic poles of the auxiliary pole type head, and the application of the system is limited. In other words, there is a disadvantage in that it is difficult to employ the auxiliary pole type head in a recording and reproducing apparatus such as a video tape recorder and a hard disc player. On the other hand, in the case of the single-sided main pole type head, the disadvantage of the auxiliary pole type head is eliminated, however, the construction of the single-sided main pole type head becomes complex. Furthermore, the single-sided main pole type head is disadvantageous in that the magnetic field distribution cannot be directed in the perpendicular direction as much as the auxiliary pole type head."}
{"text": "Cognitive decline and deficits are associated with a broad range of disorders, diseases and conditions, including decline related to normal aging, specific psychopathologies, and neurological disorders. Some cognitive deficits are related to processing of auditory and visual data, for example, and the command of motor functions while in an environment with challenging (noisy, time-limited, attentionally-demanding) conditions.\nOne deficit that distinguishes a neurotypical brain from a brain suffering from cognitive decline or a cognitive disease or disorder is the brain's ability to maintain cognitive performance in the presence of interrupting and distracting stimuli. Other deficits can involve the inability to multi-task and concentrate on performing a task in the presence of distractions.\nPrior solutions have attempted to provide a variety of computer-implemented methods, systems and tools for improving cognitive ability in aging individuals, individuals suffering from cognitive impairment, or healthy individuals wishing to enhance their cognitive abilities according to a variety of methodologies. One such methodology is interference training, which seeks to employ various tasks and stimuli to improve an individual's ability to multi-task and concentrate on performing a task in the presence of distractions. For example, U.S. Pat. No. 9,940,844B2 to Gazzaley (“Gazzaley”) provides for computer-implemented methods and tools for enhancing cognition in an individual by improving the individual's ability to deal with interference in cognitive function. Gazzaley provides for certain cognitive training methods involving presenting to an individual a task to be performed, presenting to the individual an interference, and receiving inputs from the individual. Where the interference is a distraction, the individual is to ignore the interference. Where the interference is an interrupter, the individual is instructed to respond to the interrupter as a secondary task, at which point the individual is said to be multi-tasking. Inputs are also received from the individual pertaining to this secondary task and include conducting an analysis of and/or generating feedback from the individual. The analysis includes a comparison of the performances with or without each type of interference, and, in some instances, the difficulty of the task is modulated as a function of this comparison.\nPrevious techniques for implementing interference training, such as those taught by Gazzaley, have depended on an audio-visual experience to drive a constant task and a target discrimination task. Such methods and systems are ill-suited for people who have physical impairments or difficulties, particularly visual impairment, or otherwise are subject to situations or conditions where they cannot effectively engage with a visual display.\nThrough applied effort, ingenuity, and innovation, Applicant has identified a number of deficiencies and problems with current systems, tools and methods for implementing interference training for improving one or more cognitive abilities in an individual. Applicant has developed a solution that is embodied by the present invention, embodiments of which are described in detail below."}
{"text": "Athletes engage in strenuous training to accomplish the goals of their sport. This strenuous training essentially amounts to trauma to the body, in that the human body interprets every strenuous work-out as a threat to its survival. It is known that muscle damage, caused by training, releases the catabolic hormone prostaglandin-E2. Training also causes the release of adrenocorticotropin (ACTH), which is a pituitary hormone. The presence of increased levels of ACTH increases the production of the catabolic hormone cortisol. Cortisol is also known as hydrocortisone, which is a glucocorticoid of the adrenal cortex that is a derivative of cortisone and is used in the treatment of rheumatoid arthritis. Thus, cortisol is a naturally occurring anti-inflammatory steroid. This catabolic hormone results in the release of amino acids from muscle tissue and prevents absorption of glucose. Cortisol, as a catabolic stress hormone, cannibalizes muscle tissue. High cortisol levels also result in the breakdown of connective tissue, lowered immunity and reduced muscle RNA synthesis. While cortisol may be a detriment to the athlete, scientists have conjectured that when the human body is stressed or traumatized, it triggers a \"fight or flight\" survival response. The biological design of cortisol is such that when a human is threatened, cortisol levels rise and mobilize the body for action by breaking down fat and muscle stores for emergency energy. Cortisol also reduces swelling in the event of injury. After the threat or trauma has subsided, cortisol levels return to normal. The cortisol-stress relationship is designed for intermittent physical threats and not the constant stimulation provided by today's aggressive athletes. Ongoing training results in cortisol levels that do not return to normal for extended periods of time and thereby result in the breakdown or loss of muscle tissue.\nAfter strenuous exercise, muscle tissue enters a stage of rapid nitrogen absorption in the form of amino acids and small peptides in order to rebuild the muscle fibers, grow and add new muscle fibers. During this period of repair and growth, it is important that the muscle cells have available to them sufficient levels of nitrogen in the form of amino acids. While the total level of amino nitrogen is important, the ratios of the various amino acids is even more important.\nAthletes that over-train sometimes enter into a catabolic condition. Muscle catabolism occurs when the athlete enters a negative nitrogen balance. People on diets usually have a negative nitrogen balance and therefore lose muscle when they lose weight. In contrast, a positive nitrogen balance means the animal has enough nitrogen left over to synthesize muscle proteins.\nVarious organizations have propounded a list of essential amino acids which are required on a daily basis for proper nutrition. These amino acid requirements vary throughout the growth cycle of all animals. Human muscle tissue is made up of specific amino acids and at specific ratios. If any of these amino acids are missing or deficient, the muscles will not grow, will grow slowly and may even begin to breakdown. However, if the animal is supplied with adequate amounts of protein that contain all of the muscle amino acids, this protein or source of amino nitrogen will be able to support rapid muscle recovery and growth.\nThe amino acids leucine, isoleucine and valine are the branched chain amino acids and are necessary for a positive nitrogen balance and muscle growth. The branched chain amino acids are lost at significant levels during strenuous exercise and therefore it is critical that they be available during the anabolic state.\nThe present invention is based in part, upon the discovery that the use of pyruvate in combination with an anabolic protein composition, produces a synergistic effect in increasing the lean body mass or muscle tissue of a mammal consuming same while at the same time decreasing the deposition of fat in the body.\nAs used herein and the claims, the term \"pyruvate\" means any salt or ester of pyruvic acid. Pyruvic acid has the formula: ##STR1##\nPyruvic acid is a colorless liquid with an odor resembling that of acetic acid and has a melting point of 13.degree. C. Pyruvic acid is an intermediate in the breakdown of sugars to alcohol by yeast. The mineral salts of pyruvic acid, such as magnesium pyruvate or calcium pyruvate or mixtures thereof are useful in the present invention. Sodium pyruvate is not especially preferred as it is known that sodium is associated with various negative medical conditions such as high blood pressure, water retention and heart disease. Further, certain athletes, such as bodybuilders, desire to present a defined visual image of their body which shows muscle definition and thus, the water retention properties of the sodium salt are not beneficial. Pyruvate precursors in the form of pyruvamides or pyruvyl-amino acids are also useful in the present invention. Pyruvyl-glycine is representative of the useful pyruvyl-amino acids.\nPyruvate has a number of useful applications in medicine. Pyruvate has been described for retarding fatty deposits in livers (U.S. Pat. No. 4,158,057); for treating diabetes (U.S. Pat. No. 4,874,790); for retarding weight gain (U.S. Pat. Nos. 4,812,879, 4,548,937, and 4,351,835); to increase body protein concentrations in a mammal (U.S. Pat. No. 4,415,576); for treating cardiac patients to increase the cardiac output without accompanying increase in cardiac oxygen demand (U.S. No. Pat. 5,294,641); for extending athletic endurance (U.S. Pat. No. 4,315,835); for retarding cholesterol increase (U.S. Pat. No. 5,134,162); for inhibiting growth and spread of malignancy and retarding DNA breaks (U.S. application Ser. No. 08/194,857 filed February 14, 1994); and for inhibiting generation of free radicals (U.S. application Ser. No. 08/286,946 filed Aug. 8, 1994). All of these references are incorporated herein by reference.\nU.S. Pat. 4,981,687 discloses compositions and methods for achieving improved physiological response to exercise. More specifically, this patent teaches a beverage comprising water, sugar, electrolytes, glycerol and pyruvate; and its use to ameliorate the effects of physical exertion. The teachings of U.S. Pat. No. 4,981,687 are incorporated herein by reference.\nU.S. Pat. No. 5,089,477 discloses the use of pyruvate in a liquid composition that is used to prevent weight loss in agricultural animals resulting from dehydration.\nU.S. Pat. No. 5,147,650 and U.S. Pat. No. 5,238,684 discloses and claims a fluid composition comprising water, electrolytes, sugar, glycerol, lactate and pyruvate. The teachings of U.S. Pat. Nos. 5,147,650 and 5,238,684 are incorporated herein by reference.\nU.S. Pat. No. 5,236,712 discloses and claims a beverage containing water, electrolytes, pyruvate and alanine in a concentration of from about 0.5% to about 10%. The teachings of U.S. Pat. No. 5,236,712 are incorporated herein by reference.\nPyruvate in various forms has been proposed for enteral administration and for parenteral administration. Typically, pyruvates are available in the form of salts, for example, calcium pyruvate and sodium pyruvate. U.S. Pat. Nos. 5,283,260 and 5,256,697 disclose uses for the pyruvyl-amino acids and methods for their production.\nPyruvate has been administered to mammals enterally or parenterally typically at superphysiological levels. The amount of pyruvate administered generally ranges from 1 to 20% of the mammal's caloric intake. For enteral dosage, the pyruvate may be disbursed or dissolved in a beverage product or may be included in cookies, candies or other foods. Pyruvate may also be introduced as an aqueous solution parenterally.\nThe amino acid composition of numerous proteins, biologically active polypeptides and foods, has been known for some time. One pre-eminent monograph on this subject is the Amino Acid Handbook by Richard J. Block, C. C. Thomas, publisher, 1956 (Library of Congress Catalog Card No. 56-9104). On page 343 at Table II, the approximate amino acid composition of mammalian muscle proteins on average, is set forth as follows:\n______________________________________ Amino Acid Grams/100 gm of Protein ______________________________________ Arginine 6.6 Histidine 2.8 Lysine 8.5 Tyrosine 3.1 Tryptophan 1.1 Phenylalanine 4.5 Cystine 1.4 Methionine 2.5 Serine 5.1 Threonine 4.6 Leucine 8.0 Isoleucine 4.7 Valine 5.5 Glutamic Acid 14.6 Aspartic Acid 8.0 Glycine 5.0 Alanine 6.5 Proline 5.0 Hydroxyproline 4.7 ______________________________________\nOn pages 272-273, the amino acid profile of human muscle tissue is set forth.\nProtein supplementation for serious athletes, such as body builders, is well accepted. Typically recommended dosages range between 2.0 and 3.5 gms of quality protein per kilogram of body weight per day. Numerous sources for the protein supplements are known, such as milk, egg, soy, beef and vegetable protein. Isolated fractions of these sources are also known such as ion-exchange whey protein, caseinates, whey protein concentrates, immunoglobulin and egg albumin. Protein supplements typically are provided as powders or tablets. It is also known to provide protein supplementation in the form of peptides (hydrolyzed protein) or even free amino acids. These approaches have two major limitations; cost and taste. While the use of pyruvate is known for various medical indications and the use of protein supplements are known to increase muscle mass, the prior art has failed to recognize or even consider the combination of pyruvate with an anabolic protein composition. The anabolic protein composition of the present invention is prepared through the combination of various proteins, peptides or amino acid sources to arrive at an amino acid profile that parallels the amino acid profile of human muscle tissue. Consumption of the anabolic protein composition will result in an increase in muscle mass, less catabolism of muscle tissue after strenuous exercise or as a result of disease, and quicker muscle cell restoration. Even more surprising are the synergistic results that can be achieved when pyruvate is combined with the anabolic protein composition of this invention. An additional aspect of the present invention is that the anabolic protein composition alone or in combination with pyruvate can lessen the catabolic effects of diseases such as cancer and AIDS. The prior art has not disclosed or suggested the compositions and methods of the present invention."}
{"text": "Ferrite particles are used in various applications. For example, Patent Document 1 (Japanese Patent Laid-Open No. 6-313176) discloses a fluid catalytic cracking catalyst composed of zeolite as active component and a support matrix wherein the matrix contains ferrite particles having a saturation magnetization of 10 emu/g or less and an average particle diameter of 0.001 to 20 μm.\nPatent Document 2 (Japanese Patent Laid-Open No. 2004-189513) discloses a manufacturing method of a functional ferrite particles including the steps of attaching glass particles on the surface of ferrite particles to provide a coating layer of molten glass, applying a coating of titanium oxide on the surface of the ferrite particles provided with glass followed by heat-treating. The functional ferrite particles are useful as photocatalytic material. Examples of the ferrite particles used include Mn—Zn type ferrite, Mg—Cu—Zn type ferrite, and Ni—Cu—Zn type ferrite.\nAlthough Patent Document 1 and 2 disclose that the ferrite particles are used as a catalyst, various properties of the ferrite particles are not focused.\nPatent Document 3 (Japanese Patent Laid-Open No. 2007-320847) discloses products containing a plurality of a core fine particle structure containing a plurality of primary fine particles and a plurality of primary pores, and core-shell ceramic fine particles containing a shell surrounding at least a part of the core fine particle structure. The products include a membrane, a sensor, an electrode, and a getter.\nThe core-shell ceramic fine particles described in Patent Document 3 include a core of yttrium stabilized zirconia and a shell of lanthanum ferrite. So, the ferrite particles are not small in the apparent density, various properties are not maintained in the controllable state and do not fill a specified volume with a small weight because lanthanum ferrite is used as the shell."}
{"text": "Semi-solid metal (SSM) processing is a technology that resulted from research in the early 1970's at the Massachusetts Institute of Technology. It was found that imposing a shear on a liquid metal before the solidification process began and continuing the shear while the liquid cooled below its liquidus resulted in a non-dendritic microstructure with a shear stress (and corresponding viscosity) nearly three orders of magnitude lower than that of the dendritic material. At rest, the non-dendritic metal slurry behaved as a rigid material in the two-phase region; that is, its viscosity was high enough that it could be handled as a solid. However, when a shear stress was applied, the viscosity decreased dramatically so that the material behaved more like a liquid. Thus, the slurry could flow in a laminar fashion, with a stable flow front, as opposed to the turbulent flow characteristic of molten metal.\nA property of semi-solid metal (“slurry”) that renders it superior to conventional casting processes is the non-turbulent (“laminar” or “thixotropic”) flow behavior that results when one enters the “two-phase” field of solid plus liquid. Specifically, shearing of semi-solid slurry leads to a marked decrease in viscosity, so that a partially frozen alloy can be made to flow like a non-Newtonian fluid. Thixotropic flow behavior arises from the ideal SSM microstructure of small, spherical particles (e.g., α-Al) suspended in a liquid matrix. In all semi-solid processes, it is imperative that this microstructure be produced consistently. Moreover, a uniform distribution of this microstructure throughout a volume of slurry is essential for production of high-quality components.\nThe benefits that semi-solid processing holds over conventional liquid metal casting result from the flow behavior of the partially solidified metal. The way in which a metal fills a mold (or die cavity) directly impacts the solidification of the metal; thus, the properties of the formed part can be enhanced with improved mold filling. Turbulent flow of liquid metal into a die or mold can lead to incorporation of air and mold gases into the melt. This in turn can lead to both macro- and micro-porosity in the final part, which negatively affect its mechanical properties.\nThere are several reasons that the laminar flow of semi-solid slurries is very advantageous from a casting standpoint. The first major reason is the elimination of gas entrapment, resulting in decreased porosity and oxide content in the formed part. Secondly, since semi-solid metal has lower heat content than superheated molten metal, there is less solidification shrinkage in the casting. Thus, molds can be filled more effectively and uniformly, and less post-casting machining is required. As a result, all semi-solid processes are potentially “near net-shape” processes. The reduced heat content also lowers the thermal stresses of the casting apparatus (typically a steel die) that contacts the metal, leading to longer tool life. Also, since the starting material has the thixotropic microstructure, the microstructure of any part formed with semi-solid processing is always equiaxed and non-dendritic. Therefore, the mechanical properties of the final component are better than a similar part formed from a conventional casting process.\nThe net result of the above-described advantages is that semi-solid casting can be used to produce intricate components with superior mechanical properties. The typical defects associated with molten metal casting can be circumvented when the microstructure (and thus the flow behavior) of the slurry is controlled. From an economic standpoint, it is expected that due to improved tool life, shorter cycle times, reduced machining, and ability to use less expensive heat treatment schedules, semi-solid processes will ultimately become as cost-effective as conventional casting routes such as high pressure die casting. Perhaps the most attractive attribute of semi-solid forming, however, is that due to the laminar flow of the slurry, very complex shapes can be cast, with thin walled sections on the order of millimeters.\nA number of processes have been designed to take advantage of the unique behavior of semi-solid metal slurries. These processes produce the thixotropic microstructure through some method of vigorous agitation during solidification. It was hypothesized that the induced agitation broke up (or facilitated the melting off of) dendrite arms, which then ripened and spheroidized to form a non-dendritic structure.\nThere are two routes for processing semi-solid metal, i.e. two different ways to arrive at the desired point within the solid-liquid, two-phase region. The first route starts from the solid state (“thixocasting”), and the second starts from the liquid state (“rheocasting”).\nThixocasting processes start out with a solid precursor material (“feedstock”) that has been specially prepared by a billet manufacturer, and then supplied to the casting facility. The feedstock metal has an equiaxed, non-dendritic microstructure. Small amounts or “slugs” of this alloy are partially melted by reheating into the semi-solid temperature range, leading to the thixotropic structure. In most applications, the slug is subsequently placed directly into a shot sleeve of a die casting apparatus, and the part is formed.\nDuring the initial years of SSM process development, mechanical stirring was used in various ways to break up dendrites and produce thixotropic metal structures. The combination of rapid heat extraction and vigorous melt agitation was effected by using different sizes, shapes, and velocities of stirring rods. Various researchers addressed the evolution of the “stircast” structure during this time. Although these methods worked well in that they effectively produced the desired metal structures, erosion of the stirrer became the “weak link” of the process.\nMagnetohydrodynamic (MHD) casting process has been utilized to overcome the limitations associated with the use of stirrers. In this approach, dendrites are still formed and then broken from the nuclei by agitation. The source of the agitation is not a mechanical stirrer, but alternating electromagnetic fields. Induction coils are placed around a crucible to induce these forces. The crucible is equipped with a cooling system to initiate freezing in the alloy while the melt is exposed to the electromagnetic forces. Upon cooling down to ambient temperature, the alloy has an equiaxed, non-dendritic microstructure. However, the MHD stirring process requires complicated and expensive machinery.\nThixoforming processes comprise the majority of industrial semi-solid applications used today. Rather than producing a semi-solid slurry directly from a superheated melt, a specially prepared feedstock metal is heated to form the semi-solid slurry. This approach eliminates the need for melting equipment within the SSM casting facility. However, the special feedstock must be purchased from special manufacturers at a premium in the form of metal billets, therefore thixocasting processes are not economical compared to conventional processes. Furthermore, in thixocasting processes, scrap metal must be sent back to the billet manufacturer and cannot be recycled. Most importantly, process control is difficult in thixocasting, because solid fraction (and corresponding viscosity) is sensitive to temperature gradients in the reheated material. Thus, narrow temperature ranges must be achieved consistently for successful operations. This, combined with the time it takes (several minutes on average) to reheat the feedstock to the desired solid fraction, negatively affects productivity.\nThe development of ideal one-step rheocasting applications is highly preferable to the current two- or three-step applications associated with most thixocasting methods. Current thixocasting approaches are inherently batch processes, in which only small amounts of slurry can be produced during each forming operation. This places limits on the sizes and shapes of parts produced in this manner. A continuous process would circumvent these hindrances, and could be used for a broader variety of applications.\nTo date, none of these processes has satisfactorily addressed the need for providing a continuous semi-solid casting route. The current need in the SSM field is a relatively simple, easy-to-implement, flexible process that can be used for a wide variety of processing applications. Such a process should use relatively simple methods of melt agitation to avoid the problems associated with the previously discussed approaches."}
{"text": "The present invention relates to a method for producing legible identifying data in the label portion of disc-shaped information carriers which are provided with a groove spiral for the storage of scannable signals, the method being particularly applicable to discs formed in a preliminary stage in the manufacture of records, e.g. discs used to produce record stampers.\nInformation carriers are understood to be, for example, lacquer foils, such as matrices, which serve to duplicate recorded information, by means of a stamping or pressing process, i.e. in the manufacture of audio or video records.\nRecords are legibly or visually identified by printed paper labels at their center. These labels, in addition to containing a text information about the contents of the side of the record, also contain numerical identifications, such as stock numbers and license information. In the conventional manufacturing process for records, these labels are inserted into the pressing mold before the record is pressed and are irremovably sealed with the record material during the pressing process.\nThis process stresses the paper both thermally and mechanically; it must thus meet unusual requirements. It must contain only a few volatile components, in the paper and printing dye, the printing dyes must not come off, and the paper must have a particularly high tensile strength so that it will not tear during the pressing process due to shear forces in the record material.\nIt is furthermore necessary, when a record receives labels on both sides, to apply the labels with the paper fibers running in the same direction because otherwise a bending moment results which warps the record. In exceptional cases it is also the custom to effect the labelling by glueing the labels onto the finished record after the pressing process is completed.\nA subsequent direct printing of the record surface in the label portion by means of a screen printing process is known particularly for foil-shaped pressed records. These thin plastic foils, which are generally less than 0.3 mm thick, would principally not be able to withstand, without warping, labelling with paper labels, particularly a one-sided labelling. Thus allowances must be made for the time-consuming screen printing process and the danger of subsequent damage or soiling of the finished foil by the dyes.\nFinally, it is also known to effect the numerical identification during the galvanic fabrication of the matrices, mainly so as to legibly identify the galvanically made objects. This identification consists of a catalog number constituted by a sequence of numbers and letters. The matrix contains a negative representation of this identification number and the number is then transferred to the record mass during pressing. This number cannot be read too easily on the record, but serves as a distinguishing feature to the person skilled in the art. Thus this same number is also printed on the label, for example, and provides an opportunity for checking the musical content of the record against the text of the label.\nThe drawbacks of the known identification techniques are summarized as follows:\n1. Paper labels can be used only with thicker records, not with foils of less than 0.3 mm thickness.\n2. Paper labels must be manually or mechanically placed into the pressing mold with great precision.\n3. The subsequent glueing constitutes an additional process step which also brings with it the danger of subsequent damage or soiling of the record.\n4. Screen printing is an additional process step requiring a relatively long drying period and involves the danger of soiling of the finished audio grooves by dyes.\n5. Engraved letters applied to the galvanically produced items in the conventional technique are difficult to read in the stamped surface and, due to the depth of the engraving, it would interfere with the movement of foil-type carriers during playback."}
{"text": "1. Technical Field\nThe present invention relates in general to high frequency digital systems and, in particular, to controlling the flow of information between a producer and a consumer in a high frequency digital system. Still more particularly, the present invention relates to a method and system for controlling the flow of information in a high frequency digital system from multiple memoriless producers to one or more buffering consumers via a shared resource.\n2. Description of the Related Art\nAs processor micro-architectures are optimized toward implementations that support higher clock frequencies, the complexity of the work that can be performed within each clock cycle decreases. This phenomenon has a tendency to introduce increased cycletime latency into mechanisms which control the flow of information between components in the micro-architecture, effectively delaying critical feedback within such mechanisms and eroding the aggregate bandwidth of the information flow.\nThis bandwidth erosion occurs when the buffering capacity of downstream components approaches full occupancy, such that the fill time of the remaining available capacity approaches the latency of feedback in the information flow control mechanism. In such cases, the information flow control mechanism often must pessimistically assume worst case information flow to avoid exceeding the buffering capacity of the downstream components. Thus, for example, the information flow control mechanism may assume that information is being sent (and stall information flow appropriately based on downstream capacity considerations) whether or not information is actually sent in order to avoid overflowing the downstream buffers.\nIn many cases, such bandwidth erosion can be averted by increasing the buffering capacity of downstream components to minimize the likelihood that high utilizations will dictate precautionary stalls. However, an increase in buffering capacity also raises component cost without providing additional benefit beyond a reduction in bandwidth erosion. In other cases, more sophisticated flow protocols can be introduced, which tolerate information loss due to aggressive capacity speculation by incorporating retry mechanisms. However, such mechanisms can greatly increase complexity, thus inhibiting design verification and testability. Moreover, retry protocols also increase the utilization of the upstream component, possibly requiring an increase in its size and cost.\nThe present invention offers an attractive alternative for addressing the obstacles to efficient information flow that arise in high frequency digital systems that control, monitor, or perform transformations on streaming information.\nIn accordance with the present invention, an information handling system, such as a processor, includes a plurality of producers that output packets of information, at least one consumer of the packets, and an information pipeline coupling the consumer and at least a particular producer among the plurality of producers. The information pipeline includes a shared resource having a bandwidth shared by multiple of the plurality of producers. The information handling system further includes a control unit that regulates packet output of the particular producer and that receives as inputs a producer output indication indicating that the particular producer output a packet and a shared resource input indication indicating that a packet output by the particular producer has been accepted by the shared resource for transmission to the consumer. Based upon these inputs, a number of grant messages output to the particular producer within a feedback latency of the control unit, and a portion of the bandwidth allocated to the particular producer, the control unit whether the particular producer can output a packet without packet loss. In response to a determination that the particular producer can output a packet without packet loss, the control unit outputs a grant message to the particular producer indicating that the particular producer is permitted to output a packet.\nAll objects, features, and advantages of the present invention will become apparent in the following detailed written description."}
{"text": "The present invention relates to a new and distinct cultivar of Sutera plant botanically known as Sutera hybrida and hereinafter referred to by the cultivar name ‘Balablim’.\nThe new cultivar originated in a controlled breeding program during June 2001 at Arroyo Grande, Calif. The objective of the breeding program was the development of Sutera cultivars with a trailing habit and large flowers.\nThe new cultivar was the result of the open-pollination of ‘Yasflos’, U.S. Plant Pat. No. 12,616. The new cultivar was discovered and selected by the inventor as a single flowering plant within the progeny of the above stated open-pollination during August 2001 at Arroyo Grande, Calif.\nAsexual reproduction of the new cultivar by terminal stem cuttings since 2001 at West Chicago, Ill., has demonstrated that the new cultivar reproduces true to type with all the characteristics, as herein described, firmly fixed and retained through successive generations of such asexual propagation."}
{"text": "Various medicaments have been used in the form of inhalation formulation for the treatment of respiratory diseases, e.g., asthma and chronic obstructive pulmonary disease (COPD). A particular advantage of inhalation formulation is that only a small amount of a pharmaceutically active ingredient is required to achieve the desired therapeutic effect; however, there are drawbacks to the formulation that only a part of the pharmaceutically active ingredient administered will be delivered to a target site, or there is a great possibility that the pharmaceutically active ingredient will be delivered to sites where no treatment is required, thereby causing adverse side effects. Thus, continuous efforts are being made to maximize the therapeutic effect of the formulation so as to achieve reliable targeted delivery to the site where the therapeutic effect is desired and to prevent the delivery of the pharmaceutically active ingredient to the site where no treatment is required.\nFor effective administration of inhalation formulation, inhalers, which administer the drug by sucking in the air with the drug and delivering them into the air passage, have been widely used for treatment of respiratory diseases. The most common inhaler systems are metered dose inhalers (MDI), which had been used extensively since its approval in 1956, and once occupied 80% of the inhaler market. However, a rise of environmental concerns, e.g., depletion of ozone layer and global warming, has shifted research interests to focus on dry powder inhalers (DPI) in recent years. In the current stage, researchers are concentrating their efforts to remedy the shortcomings of MDI formulations by employing DPI formulations. MDIs typically comprise pharmaceutically active ingredients and a solvent as a propellant in a compressed state, which deteriorates its stability; and the spraying speed is fast, and thus it reaches the laryngopharyngeal space too fast. DPIs, however, are easy to use; and only comprised of powder solid particles, and therefore are advantageous in terms of stability (see Martin J Telko and Anthony J Hickey, Dry Powder Inhaler Formulation, Respiratory Care, September 2005, Vol 50, No. 9).\nMeanwhile, various drugs are being tested for the prevention and treatment of respiratory diseases. For example, a selective beta-2 adrenoceptor agonist (beta-2 agonist) can induce bronchodilation, and can be used to relieve respiratory distress. Beta-2 agonists may be broadly divided into short-acting beta-2 agonists and long-acting beta-2 agonists. Short-acting beta-2 agonists, e.g., salbutamol, fenoterol, levalbuterol, terbutaline, etc., provide immediate relief, but their reaction time is rather short. In contrast, long-acting beta-2 agonists, e.g., formoterol, indacaterol, salmeterol, tulobuterol, etc., provide sustained bronchodilation, but patients are required to take them two or more times per day because the normal reaction time of these drugs is less than 12 hours.\nBeta-2 agonists can alleviate bronchoconstriction in patients, but other drugs, e.g., steroids, are used to treat inflammation, which is another cause of asthma. Examples of steroids include inhaled corticosteroids (ICS) such as beclomethasone, budesonide, flunisolide, fluticasone propionate, mometasone furoate, triamcinolone, and the like.\nAlso, another type of drug called an inhaled anticholinergic is well-known as a stable and effective bronchodilator which can be used for treatment of COPD. Anticholinergic agents can increase the level of forced expiratory volume in 1 second (FEV1), prevent static or dynamic hyperinflation (overexpanded lung), and reduces exacerbations of COPD. There is a limited number of inhaled anticholinergic bronchodilators that are currently available, e.g., rapid-onset types such as ipratropium bromide, oxitropium bromide, etc., and long-acting types such as tiotropium bromide, etc.\nGlobal Initiative for Asthma (GINA) and Global Initiative for Chronic Obstructive Lung Disease (GOLD) are suggesting an incremental treatment method based on the progression of the disease condition, which includes the use of a combination formulation of drugs having different or complementary action mechanisms. For instance, a long-acting beta-2 agonist is prescribed to patients with asthma or COPD having FEV1 level of less than 80%, and COPD patients with accompanying respiratory distress having FEV1 level of less than 50% or who are experiencing frequent acute exacerbations are prescribed with ICS in addition to beta-2 agonists.\nA number of combinations of the aforementioned drugs are known already, and one typical example is an inhalation formulation comprising salmeterol xinafoate and fluticasone propionate (Seretide, GSK). Currently, Seretide is available in MDI (Evohaler) and DPI (Diskus) formulations. Seretide provides an effective bronchodilation induced by the long-acting beta-2 agonist salmeterol, as well as potent anti-inflammatory action caused by the ICS, fluticasone propionate. In the case of Seretide Diskus formulation, which is provided in the form of DPI formulation, both beta-2 agonist and inhaled corticosteroid may be inhaled at once, but the formulation does not show sustained bronchodilation action, and thus patients are required to take the formulation two or more times per day. Another drawback of this formulation lies in that the amount of the excipient is too small to give a sensation in the lungs upon administration, and sometimes the dose is not properly delivered or is taken two or more times because it is impossible to observe administered formulation.\nAlso, a combination therapy comprising a rapid-onset anticholinergic ipratropium bromide and a long-acting beta-2 agonist salmeterol is disclosed in WO01/76601, and an additional combination therapy using an anticholinergic, beta-2 agonist as well as steroid is disclosed in U.S. Pat. No. 6,423,298 and WO02/7672.\nNevertheless, said formulations do not relate to a triple combination formulation which can exert fast-acting bronchodilation induced by a beta-2 agonist, anti-inflammatory action by a corticosteroid, and sustained bronchodilation by an anticholinergic at once.\nRecently, KR Patent Laid-Open Publication Nos. 10-2010-0063116 and 10-2009-0121338 mentioned a triple combination formulation of beta-2 agonist, corticosteroid, and anticholinergic. However, they neither consider any specific device, effective dose amount, manufacturing method thereof, packaging type, particle size of the carrier material, etc., nor provide assessment data thereof. Pressure drop values of inhalation formulations, especially in the form of dry powder inhalation formulations, vary when different types of devices are used, and the amount of active ingredient delivered to lungs may vary with packaging forms, e.g., blister packaging vs. capsule packaging. The properties and ratio of excipient (such as lactose, etc.), which is used as a carrier, can also cause a large difference in therapeutic effects. Moreover, even if drugs from the same drug group were used, undesirable results such as deterioration in uniformity and storage stability may occur depending on physicochemical properties of the respective drugs.\nAlthough some drugs and a combination formulation thereof for the prevention or treatment of respiratory diseases are known, there are no specific compositions or preparation method thereof developed for a triple combination formulation which can administer a long-acting beta-2 agonist, an inhaled corticosteroid and an anticholinergic together at once. Thus, there has been a need to develop a composition of a composite formulation, which can stably and accurately administer said three groups of drug in a single dose, to improve patient compliance and enhance patients' convenience to carry the formulation."}
{"text": "Mass spectrometry is a very sensitive analytical method used for important research and for applications of analytical chemistry, particularly life science. Electrospray ionization (ESI) is generally regarded as the best-characterized and most efficient method for ionization of molecules in solution phase. The process can be conveniently divided into three stages: droplet formation, droplet evaporation and ion formation (Gaskell, S. J. Journal of Mass Spectrometry 1997, 32, 677-688). When a strong electric field is applied to a solution flowing through a mass spectrometer probe, a Taylor cone is formed at the tip of the probe, resulting in a mist of small droplets being emitted from the tip of this cone. Due to the evaporation of the free droplets and Coulombic forces, ions of sample analyte are produced. The ions enter a mass spectrometer and are subsequently analyzed.\nA problem with ESI is that sample preparation is still a necessary step before ESI can be used for analysis of many types of samples. Prior to analyzing a sample by ESI mass spectrometry, the sample will undergo extraction and filtration protocols to purify the sample, for example to remove salts and detergents. Such protocols are complex, time-consuming, and expensive. Further, reagents used during the purification process can interfere with subsequent analysis of a target analyte in the purified sample. Additionally, samples that are not in solution must be dissolved as well as purified prior to ESI analysis.\nMore recently, the concept of ambient ionization has been developed, and now this family of ambient ionization has more than twenty members, such as desorption electrospray ionization (DESI) and direct analysis in real time (DART). Ambient ionization by mass spectrometry allows the ionization of analytes under an ambient environment from condensed-phase samples without much or even any sample preparation and/or pre-separation, offering a solution for real time and in situ analysis for complex mixtures and biological samples. These ambient ionization methods are leading are extending the mass spectrometry revolution in life science, environment monitoring, forensic applications and therapeutic analysis. However, the above described ambient ionization techniques still require pneumatic assistance, a continuous flow of solvent, and a high voltage power supply for the analysis of samples.\nThere is an unmet need for systems and methods that can combine sample preparation and pre-treatment and the ionization process for mass analysis of samples that do not require pneumatic assistance or a continuous flow of solvent for the analysis of the samples."}
{"text": "The present Invention relates to a new and distinct cultivar of Poinsettia plant, botanically known as Euphorbia pulcherrima Willd., and hereinafter referred to by the name ‘PER7102’.\nThe new Poinsettia is a naturally-occurring whole plant mutation of the Euphorbia pulcherrima Willd. cultivar Eckakeem, disclosed in U.S. Plant Pat. No. 13,838. The cultivar PER7102 was discovered and selected by the Inventor as a flowering plant within a population of plants of the parent cultivar in a controlled environment in Encinitas, Calif., on Jan. 7, 2002. The new Poinsettia was selected on the basis of its flower bract coloration.\nAsexual reproduction of the new Poinsettia by terminal cuttings propagated in a controlled environment in Encinitas, Calif., since March, 2002, has shown that the unique features of this new Poinsettia are stable and reproduced true to type in successive generations of asexual reproduction."}
{"text": "The invention relates to a camshaft adjusting mechanism for adjusting the phase position of a camshaft relative to a camshaft drive or crankshaft of an internal combustion engine.\nVariable valve timing in order to optimize the fuel consumption and engine power output in motor vehicles is already known. It is possible, for example, to adjust the phase position of a camshaft in relation to the crankshaft in order to vary the valve actuation. Such arrangements for adjusting the phase position usually take up a relatively large overall space, which is not readily available in modern vehicles accommodating numerous assemblies. In order to adjust the phase position, use may be made, for example, of a hydraulic vane adjuster, the rate and angle of adjustment of which, however, are very limited. At low engine speeds and low oil pressure, and when the vehicle is at standstill, the phase position cannot be adjusted.\nThese problems can be avoided by using electrically driven adjusting mechanisms. The published patent application DE 100 38 354 A1 discloses an electrically driven adjusting mechanism having an adjusting device for adjusting the phase position of the camshaft in relation to a crankshaft driven by an internal combustion engine, together with a drive gear, driven by the crankshaft, for driving the camshaft, in which a swash plate mechanism is provided. Such an arrangement, however, entails high costs because of the highly geared two-stage transmission and the servomotor which are required.\nDE 41 33 408 A1 discloses a camshaft adjusting mechanism in the form of a four-shaft transmission in which two drive shafts are provided for adjustment, which can be simultaneously subjected to a brake torque.\nIt is the object of the present invention is to provide a very compact and inexpensive camshaft adjusting mechanism."}
{"text": "The present invention relates to an arbor for lathe cutting lenses, especially contact lenses, from a blank or button, and a method and apparatus for making the arbor.\nA conventional method of manufacturing lenses, especially contact lenses, involves lathing the lens from a cylindrical blank of polymerized lens material (such cylindrical blanks commonly referred to as a xe2x80x9cbuttonxe2x80x9d). The buttons may be cut initially from rod stock or sheets, or individually cast in cups using a curable liquid monomeric composition. Additionally, lenses may be lathe cut from a semi-finished lens blank, i.e., a blank having one molded-in lens surface, where the opposite lens surface is lathed to obtain a second desired lens surface and an article with a final contact lens shape.\nWhile it is possible in some cases to insert the lens blank directly into the lathe collet, it is more typical to first attach the button to a separate pin or xe2x80x9cblockxe2x80x9d with an adhesive, the opposite end of the block being configured for removable insertion into the collet of the lathe. While the block and button are turned, the desired concave or convex curve is lathed into the exposed end of the blank. Eventually, the lathed article must be separated from the block (a process often referred to as xe2x80x9cdeblockingxe2x80x9d). However, a drawback of such lathe cutting operations is that the blocking and deblocking steps may consume a significant amount of time and are prone to error, especially when it is necessary to maintain axial alignment between the block and the article being lathe cut.\nU.S. Pat. No. 4,921,205 discloses a process where a semi-finished blank is cast between a plastic mold cup and a radial die, for example, the radial die may form a finished concave lens surface. The radial die is removed, with the cast blank remaining adhered to the mold cup. The mold cup and blank assembly may be shipped to a lab for manufacture of a contact lens, by first machining the mold cup off the blank, then machining away the extraneous portion of the blank to form the second, opposite lens surface. Presumably, the cast surface of the blank would be adhered to a block for these machining operations, similar to other conventional methods for lathe cutting a lens surface from a semi-finished blank.\nU.S. Pat. No. 5,110,278 discloses a process for producing toric contact lenses. A semi-finished blank is cast in an arbor having a cup-shaped top portion. The arbor also includes a hollow bottom portion integrally formed as a one-piece molding with the top portion. This hollow bottom portion is mounted on a lathe spud, for lathe cutting the exposed front surface of the semi-finished blank retained in the cup-shaped top portion. Each arbor is injection molded, and the arbors are not reusable since the lathing operation involves simultaneous cutting of the lens blank and sidewalls of the casting cup top portion.\nThe present invention provides an arbor for lathe cutting lenses, especially contact lenses, from a blank or button, and a method and apparatus for making the arbor. The arbor is designed for the mounting of a lens mold with a cast lens blank therein during lathing of the lens blank, and the arbor is made in a manner to ensure the lens mold is axially aligned with the arbor. The arbor can be reused for lathe cutting of multiple lens blanks retained in their respective lens molds.\nAccording to a first embodiment, this invention provides a method for making an arbor comprising: providing a lens mold that includes a lens-forming molding surface and an opposed non-molding surface, and placing a hardenable material (for example, liquid or pliable solid material) between the opposed surface of the lens mold and an end of an arbor; moving the lens mold and the arbor towards one another while maintaining axial alignment therebetween, whereby the hardenable material is pressed between and contacts a surface of the lens mold and a surface of the arbor end and conforms to said surfaces; hardening the material to form a hardened replica of the lens mold surface contacted by said material; and separating the arbor with the hardened material attached thereto from the lens mold. The arbor may be placed on a lathe collet, a second contact lens mold including a lens blank cast therein may be secured to the hardened material of the arbor, and a desired lens surface may be lathe cut in the lens blank. This same hardened material of the arbor may be used repeatedly to lathe cut a desired lens surface in multiple lens blanks retained in their respective lens molds.\nAccording to preferred embodiments, the central axis of the lens mold is congruent with the central axis of the arbor when said material conforms to the lens mold and arbor end surfaces. Also, the lens mold may be secured in a first fixture to inhibit axial movement thereof, and the arbor may be secured in a second fixture to inhibit axial movement thereof. These first and second fixtures may include complementary tapered surfaces for maintaining axial alignment of the lens mold and arbor secured thereon. The first and second fixtures may include a bore for securing the lens mold and arbor, respectively.\nAccording to other preferred embodiments, the lens mold includes a sidewall depending from the opposed, non-molding surface, such that the hardenable material contacts and conforms to an inner surface of the sidewall, and wherein the arbor end is received in a space surrounded by the sidewall. At least a portion of the sidewall inner surface is preferably noncylindrical.\nAccording to a second embodiment, this invention provides a method comprising: casting a lens blank in a lens mold; securing the lens mold with the blank on an arbor end, the arbor end including a replica of an inner surface of the lens mold; and lathe cutting a lens surface in the blank.\nAccording to other embodiments, the invention provides an arbor comprising a first end for mounting on a lathe collet and an opposed second end, and a hardened material surrounding and adhered to the second arbor end, the hardened material having an outer surface that is a replica of an inner surface of a contact lens mold."}
{"text": "Field of the Invention\nThe invention relates to a method and configuration for transmitting digital data, which is divided up into data frames of variable length, from a first data bus to a second data bus and controlled by a microprocessor.\nThe integrated services digital network (ISDN) provides each subscriber with access to two so-called B channels for communicating with other subscribers, as well as to a so-called D channel in order to interchange signaling for a desired connection with the ISDN. While the ISDN does not ensure any data protection for the subscribers when communicating on the B channels, the ISDN demands for the D channel data protection conforming to the requirements of Layer 2 of the international standardization organization/open systems interconnection (ISO-OSI) layer model. To this end, the high-level data link control (HDLC) protocol is used for the ISDN.\nHence, the terminals and network interface equipment for the ISDN are provided with HDLC control devices which check the data protection information for a received D-channel frame and add protection information to a D-channel frame which is to be transmitted. The applicant markets such control devices under the product name PEB 2070 which is an integrated semiconductor component.\nAn incoming D-channel signal is supplied to an HDLC receiver logic unit, which checks the protection information for the received D-channel signal. An output signal from the HDLC receiver logic unit is supplied to a first-in first-out (FIFO) memory of a fixed size. A microprocessor removes the data contained in the FIFO memory from the latter in blocks. In the transmission direction, the microprocessor writes data to an FIFO transmission memory in blocks. An HDLC transmission logic unit removes the data contained in the FIFO transmission memory from the latter, adds protection information to a complete D-channel frame and outputs the D-channel signal thus obtained."}
{"text": "I. Field\nThe present disclosure relates generally to communication, and more specifically to techniques for supporting communication in a wireless communication system.\nII. Background\nWireless communication systems are widely deployed to provide various communication content such as voice, video, packet data, messaging, broadcast, etc. These wireless systems may be multiple-access systems capable of supporting multiple users by sharing the available system resources. Examples of such multiple-access systems include Code Division Multiple Access (CDMA) systems, Time Division Multiple Access (TDMA) systems, Frequency Division Multiple Access (FDMA) systems, Orthogonal FDMA (OFDMA) systems, and Single-Carrier FDMA (SC-FDMA) systems.\nA wireless communication system may include a number of base stations that can support communication for a number of user equipments (UEs). A UE may communicate with a base station via the downlink and uplink. The downlink (or forward link) refers to the communication link from the base station to the UE, and the uplink (or reverse link) refers to the communication link from the UE to the base station."}
{"text": "Typical properties required of steel sheets for use in domestic electrical appliances, automobiles, furniture, building and the like include strength and resistance to rusting. Parts of television cathode-ray tubes, such as explosion-proof bands and support frames, should shield the influence of geomagnetism so that electron beams, when passed through a space constituted thereby, are not deflected. Improved geomagnetic shielding properties referred to herein mean that the relative permeability in a d.c. magnetic field around 0.3 Oe corresponding to geomagnetism is large. Use of steel sheets satisfying this property requirement even in automobiles, which have more and more become electronically controlled, creates a possibility of the prevention of erroneous actuation of instruments.\nIn general, improved geomagnetic shielding properties can easily be realized by using non-oriented magnetic steel sheets such as specified in JIS C 2552. In this case, what is required is only to increase the relative permeability in a d.c. magnetic field around 0.3 Oe corresponding to geomagnetism, and, unlike rotary machines, properties in high magnetic field are not required. If such steel sheets could be produced in the same equipment as used for the production of steel sheets for press working, the thickness range of producible sheets could be broadened and, in addition, the production cost could be reduced.\nReducing fine precipitates present in steel or coarsening ferrite grains to facilitate the movement of domain walls is known to be effective in increasing the relative permeability in a d.c. magnetic field around 0.3 Oe corresponding to geomagnetism. For example, Japanese Patent Laid-Open Publication No. 61330/1991 discloses a method wherein a low carbon aluminum killed steel is subjected to open coil decarburization annealing to coarsen grains. Further, Japanese Patent Publication No. 6134/1996 and Japanese Patent Laid-Open Publication No. 27520/1996 discloses a method wherein a steel having a carbon content reduced to not more than 0.01% with reduced impurities is continuously annealed to coarsen grains. For steel sheets produced by these methods, however, the yield point is estimated to be about 250 MPa at the highest.\nOn the other hand, when reducing the amount of steel products used is contemplated from the viewpoints of a reduction in weight and life cycle assessment (LCA), a high yield point of, for example, 250 to 300 MPa or more, is required. This necessitates enhancing the yield point through utilization of one of or a combination of two or more of solid solution strengthening, fine grain strengthening, precipitation strengthening, and work strengthening. In any case, however, geomagnetic shielding properties rapidly deteriorate with increasing the yield point. Further, when the silicon content is increased, plates are likely to be broken at the time of rolling, leading to lowered productivity and yield. This makes it impossible to attain the object.\nAccordingly, it is an object of the present invention to solve the above problems of the prior art and to provide high strength cold rolled steel sheets and high strength plated steel sheets having improved geomagnetic shielding properties, that is, high relative permeability in a d.c. magnetic field around 0.3 Oe, and a process for producing the same. It is another object of the present invention to provide explosion-proof bands for television cathode-ray tubes or outer magnetic shielding materials using these high strength cold rolled steel sheets and high strength plated steel sheets. Cold rolled steel sheets and plated steel sheets refer to cold rolled steel sheets not subjected to surface treatment in the narrow sense which are used in domestic electrical appliances, including explosion-proof bands and support frames for television cathode-ray tubes, automobiles, furniture, building and other applications, and surface treated steel sheets, for example, plated steel sheets subjected to plating for rust preventive purposes, for example, electroplated steel sheets subjected to plating with zinc or zinc-nickel, galvanized steel sheets, and alloyed galvanized steel sheets, and plated steel sheets subjected to treatment for further improving press forming properties and rusting resistance, such as plated steel sheets subjected to alloying of the plating and plated steel sheets with an organic coating formed as an upper layer."}
{"text": "The present invention relates to a drive motor for rotating a multi-faceted mirror, and more particularly, to a drive motor for rotating a multi-faceted mirror, which is used to change the path of a laser beam in a laser printer or a laser scanner.\nAn optical data processing device such as a laser printer or a laser scanner has been developed to process optical data on a large scale. The optical data processing device for printing or scanning a large amount of data at high speed is comprised of a laser 1 with excellent convergence, a high-speed beam deflector 2, rotating a multi-faceted mirror 3, and a drive motor 4 for rotating the multi-faceted mirror 3, as shown in FIG. 1.\nFIG. 2 is a simplified view of a bearing portion of a conventional drive motor for rotating a multi-faceted mirror. As shown in FIG. 2, a pair of upper and lower ball bearings 40 are disposed between a bearing housing 10 and a shaft 30 fixed to a rotor 41 of the motor. The reference numeral 11 denotes a bearing cover.\nThere are no problems during printing or scanning with drive motors having ball bearings 40 installed therein, as long as the printer or scanner operates at a low speed. However, significant problems occur when the printer or scanner operates at a high speed. That is, if printing or scanning is performed by rotating the multi-faceted mirror 3 at above speeds of 20,000 RPM, there is significant shaking and noise in the ball bearings 40, which are contact-type bearings."}
{"text": "The compounds of formula (I) are used as microbicides for treating plant diseases caused by Pythiaceae or Peronosporaceae. The compounds of formula (I) were already disclosed in Korean Patent Laid-open Publication No. 94-19960 and the corresponding foreign applications, for example, U.S. Pat. application Ser. No. 08/287,917, JP Patent Application No. 192529 and EP Patent Application No. 94112652.6 which were filed by the present applicant.\nFurther, a process for preparation of 2-aminothiazolecarboxamide derivatives including the compounds of formula (I) using 2-aminothiazole carboxylic acid as an intermediate was described in Korean Patent Laid-open Publication No. 97-24120. However, this process has the disadvantage that it is not economic when applied to industrial production due to the many steps for obtaining the intermediate and the low yield.\nAccordingly, the present inventors have studied to improve the prior process by solving the above mentioned problems, and as a result, have brought the present invention to completion."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for forming an image, such as a color laser printer.\n2. Background Art\nAs an apparatus for forming a multicolor image, such as a four-color image, on a recording medium, such as recording paper, a tandem type apparatus has been known, which has, for each of colors, a photoreceptor, an exposing unit for exposing the photoreceptor to form an electrostatic latent image on the surface thereof, and a developing unit for developing the electrostatic latent image by attaching a charged developer on the surface of the photoreceptor having the electrostatic latent image formed thereon. A four-cycle type apparatus has also been known, which has one pair of an exposing unit and a photoreceptor, and plural developing unit for the respective colors arranged around the photoreceptor.\nThe four-cycle type apparatus is not suitable for speeding up the image formation because the exposing and developing steps of the photoreceptor are sequentially carried out with changing colors. In the tandem type apparatus, on the other hand, the exposing and developing steps of the photoreceptors can be simultaneously carried out for the colors, and a multicolor image can be formed by sequentially transferring the developers attached on the photoreceptors corresponding to the respective colors on a recording medium. Therefore, the tandem type apparatus for forming a multicolor image exhibits a color image forming speed equivalent to the case of a monochrome image, and is suitable for high-speed image formation."}
{"text": "1. Technical Field\nThe invention relates to cable television (CATV) interdiction systems and, more particularly, to a method and apparatus for controlling the jamming parameters in such systems.\n2. Description of the Prior Art\nIn the past, a scrambler has been provided to encode premium television channels at a headend of a cable television system. The applied scrambling precluded reception by an unauthorized converter/decoder at a connected premises. Data representing the channels or tiers of programming to which the subscriber was entitled were addressably transmitted to a particular converter/decoder and stored in an authorization memory. As a result of the addressed transmission, a subsequently transmitted program would be authorized by selectively enabling the decoder portion of the converter/decoder to decode the scrambled premium channel or program.\nThe provision of one scrambler per premium channel at the headend and the inclusion of a descrambler in each converter/decoder at the premises of the television receiver was particularly expensive. Furthermore, providing a converter/decoder on premises has turned out to be a great temptation to service pirates who imaginatively seek ways to receive premium channels. As a result, cable television equipment manufacturers have entered into a veritable war with such pirates resulting in complicated service authorization protocols, which in some instances involve multiple layers of encryption by both in-band and out-of-band data transmissions thereby further increasing the costs of the converter/decoder. In addition, scrambling systems may leave artifacts in the final signal.\nConsequently, the cable industry has reviewed other technology developed in the early stages of cable television, such as the application of negative and positive traps, and more recent techniques, such as interdiction, to improve CATV systems.\nA relatively recent technique for premium channel control is the interdiction system, so called because of the introduction of an interfering signal into a premium channel at the subscriber's location. Most embodiments consist of a pole-mounted enclosure located outside the subscriber's premises designed to serve one or more subscribers. This enclosure contains at least one microprocessor controlled oscillator and switch control electronics to secure several television channels. Control is accomplished by injecting an interfering or jamming signal into unauthorized channels from this pole-mounted enclosure.\nFor the sake of efficiency, it is known to utilize one oscillator to jam several premium television channels. This technique not only reduces the amount of hardware required, but also maximizes the system flexibility. The jamming signal frequency is moved as a function of time from channel to channel. The oscillator is frequency agile and hops from jamming one premium channel frequency to the next. Cable television channels and, of course, premium service channels may extend over a wide range of frequencies, for example, from 54 Mhz. up to including 1000 Mhz. Thus, if only one oscillator were provided, it would have to be frequency agile over a wide range.\nOne such system, illustrated in U.S. Pat. No. 4,450,481 by Dickenson, has a single frequency agile oscillator which provides a hopping gain-controlled jamming signal output to four high frequency electronic switches. In this system, each switch is associated with one subscriber drop. Under microprocessor control and depending on which subscribers are authorized to receive transmitted premium programming, the microprocessor selectively gates the jamming signal output of the single oscillator via the switches into the path of the incoming broadband television signal to each subscriber. Consequently, an unauthorized subscriber upon tuning to a premium channel will receive the premium channel on which a jamming signal of approximately the same frequency has been superimposed.\nA significantly more advantageous interdiction system is disclosed in U.S. Pat. No. 4,912,760 by West, Jr., et al. which controls a plurality of frequency agile oscillators for each subscriber. Each voltage controlled oscillator is allocated a continuous band of frequencies consistent with the elimination of jamming signal harmonics which could disturb authorized programming at a higher channel frequency. The interdiction apparatus includes generating and storing frequency control words for operating the voltage controlled oscillators consistent with a headend selected jamming factor for a particular channel to be jammed and addressably transmitted and stored premium programming authorization data. In U.S. application Ser. No. 07/476,041, filed Feb. 6, 1990, by West, Jr., et al., the method of programming the time slots of the West I system to vary the dwell time and jamming factor of the jamming signals for different premium channels is disclosed.\nThe system is West I and West II was configured in such a way that different channels can be programmed with a variable jamming factor. This provides a channel with a programmable number of time slots which can be assigned and a variable dwell time for a given hopping rate. This variability is necessary because it has been recognized that the effectiveness of a jamming signal should be different for specific channels because not all channels have the same information content. Interdiction works most effectively by destroying the entertainment value of the channel and can be accomplished in a number of ways.\nFor example, channels which have considerable video content may have to be jammed with a larger number of time slots than those which do not. This is true for many types of premium viewing including adult type channels where the video portion should be entirely obscured. However, for many other types of channels jamming with single or a small number of time slots provides enough interference with the channel information to discourage watching unless the jamming signal is is removed.\nIt has additionally been found that other types of variations in the jamming parameters of an interdiction system can increase jamming effectiveness. A pulse jamming method has proved to be effective in many instances where total obscuration of the video is not needed. In pulse jamming, the jamming signal for a specific channel is turned on and interdicts the channel for a particular duration and then is turned off for a particular duration. This pulsing of the jamming signal provides enough interference with the channel that, even when the channel is in the clear for short periods of time, its entertainment value is lost. Pulse jamming is more difficult to do in a time slot based system such as West I and West II because the jamming energy is already being time shared between several channels. To have an effective and efficient interdiction system, the jamming energy available when a channel is in the clear when operating in pulsed mode should be reallocated to other jammed channels without having to reprogram the time slots. Therefore, it would be advantageous to provide a method of combining pulse jamming with the programmable time slot jamming technique of West I and West II.\nFurther, it has been noted that the hopping rate can effect the jamming effectiveness of a system. For channels with a greater jamming factor (more time slots) it is more effective to have a faster hopping rate. For channels with a lower jamming factor (less time slots) it is more effective to have a slower hopping rate. With faster hopping rates and more time slots the obscuration of the video signal is ensured. With a slower hopping rate and less time slots, many television receivers tend to lose vertical synchronization so that the video portion of the program will roll.\nBecause of the flexibility in the channel line-up and multitude of service tiers demanded by today's subscribers, system operators need to have a combination of these methods to discourage unauthorized viewing of premium channels. It would be advantageous to provide an interdiction system with a variety of variable jamming parameters such as variable dwell time, variable hopping rate, and pulsed mode jamming so as to be able to modify jamming effectiveness based on channel line-ups, number of pay channels and other system factors. It would be further advantageous to provide combinations of these variable jamming parameters to generate the optimal jamming effectiveness for a particular system configuration."}
{"text": "During a manufacturing process of various structures, devices, systems, and the like, two separate and distinct components are secured together by one or more fasteners. The components may not have similar sizes or shapes. As an example, one of the components may be a flat planar piece, while the other component may be curved. In order to secure the two components together, a two piece washer assembly may be used. The washer assembly includes a first loose washer that receives a second loose washer. The two loose washers are configured to move through a limited range of motion in order to accommodate the angular differences between the two components.\nIn order to securely connect the two components together, an individual manipulates the washer assembly on one side of a first component, and a fastener on another side of the component. However, the washer assembly is susceptible to separating before receiving the fastener.\nAs can be appreciated, during the manufacture of a large system, such as an aircraft, the separate and distinct washers may be lost. For example, a flat drill plate may be secured to a portion of a curved fuselage panel. A first individual may be positioned on one side of the fuselage panel to secure the washer assembly, while a second individual may be positioned on an opposite side of the fuselage with a fastener. For large structures such as aircrafts, two individuals often work together to secure components (such as a drill plate and fuselage panel) together, through the use of the two piece washer assembly and a fastener. However, the unconnected loose washers may separate and fall into inaccessible areas of the fuselage. As can be appreciated, the manufacturing process may be labor and time intensive, and may be susceptible to loose parts being lost, or misplaced, such as within the system being manufactured."}
{"text": "1. Field\nThis disclosure relates generally to semiconductor devices, and more specifically, to a 3-D semiconductor die structure with containing feature and method.\n2. Related Art\nWhen forming metal to metal die to die interconnections, corresponding pads of each die are sized large enough to overlay one another, taking account of placement misalignment. For large pads that are equally sized, bonding metal like tin (Sn) can run on the surface of the die, thus undesirably shorting adjacent interconnects not meant to be shorted.\nThere exists a problem of die movement after pick and place of die onto a wafer and during simultaneous bonding of die at the wafer level. The die movement can cause connect to connect shorting or misaligned opens during thermal compression die to wafer bonding.\nAccordingly, there is a need for an improved method and apparatus for overcoming the problems in the art as discussed above."}
{"text": "Reconstruction of three-dimensional objects in a scene from multiple two-dimensional images of the scene has been the subject of research since the late 19th century. Such reconstruction may be useful in a number of areas, including obtaining information about physical (three-dimensional) characteristics of objects in the scene, such as determination of the actual three-dimensional shapes and volumes of the objects. Reconstruction has also recently become particularly important in, for example, computer vision and robotics. The geometric relation between three-dimensional objects and the images created by a simple image recorder such as a pinhole camera (that is, a camera without a lens) is a source of information to facilitate a three-dimensional reconstruction. Current practical commercial systems for object reconstruction generally rely on reconstruction from aerial photographs or from satellite images. In both cases, cameras are used which record images from two locations, whose positions relative to a scene are precisely determinable. In reconstruction from aerial photographs, two cameras are mounted with precise spacing and orientation on a common airborne platform, which ensures that the geometries of the cameras relative to each other are fixed in a known condition. With satellites, the positions and orientations of the satellites can be determined with great accuracy, thereby providing the required geometrical information required for reconstruction with corresponding precision. In any case, reconstruction of the desired objects shown in the images can be performed from two-dimensional photographic or video images taken from such an arrangement.\nGenerally, reconstruction methods are non-linear and they generally do not behave well in the presence of errors in measurement of the various camera calibration parameters and in the images from which the objects are to be reconstructed. Conventional reconstruction methods rely on the successful decoupling of two sets of parameters known as intrinsic and extrinsic parameters. The extrinsic parameters are related to the external geometry or arrangement of the cameras, including the rotation and translation between the coordinate frame of one camera in relation to the coordinate frame of the second camera. The intrinsic parameters associated with each camera is related to the camera's internal geometry in a manner that describes a transformation between a virtual camera coordinate system and the true relationship between the camera's image plane and its center of projection (COP). The intrinsic parameters can be represented by the image's aspect ratio, the skew and the location of the principal point, that is, the location of the intersection of the camera's optical axis and the image plane. (Note that the camera's focal length is related to the identified intrinsic parameters, in particular the aspect ratio, and thus it need not be considered as a parameter.)\nThese intrinsic and extrinsic parameters are coupled together and it is possible to recover the Euclidean three-dimensional structure of the scene depicted in two views only if these two sets of parameters can be decoupled. The precise manner in which the intrinsic and extrinsic parameters are coupled together is as follows. If the intrinsic parameters for the cameras are used to form respective three-by-three matrices M and M', and R and \"t\" represent the rotational and translational external parameters, then for points p=(x,y,1).sup.T and p'=(x',y',1).sup.T (\"T\" represents the matrix transpose operation) representing the projection in the two images of a single point P in the scene, EQU z'p'=zM'RM.sup.-1 p-M't (1)\nwhere z and z' represent respective depth values for point P relative to the two camera locations.\nThere are several general methods for reconstruction. In one set of methods, the values of the various parameters in equation (1) are determined. In one such method the values of the internal parameters are determined by a separate and independent \"internal camera calibration\" procedure that relies on images of specialized patterns. In a second such method, more than two views of a scene are recorded and processed and the two sets of parameters are decoupled by assuming that the internal camera parameters are fixed for all views. Processing to determine the values of the parameters proceeds using non-linear methods, such as recursive estimation, non-linear optimization techniques such as Levenberg-Marquardt iterations, and more recently projective geometry tools using the concept of \"the absolute conic.\" One significant problem with the first approach (using a separate internal camera calibration step) is that even small errors in calibration lead to significant errors in reconstruction. The methods for recovering the extrinsic parameters following the internal calibration are known to be extremely sensitive to minor errors in image measurements and require a relatively large field of view in order to behave properly. In the second approach (using more than two views of a scene) the processing techniques are iterative based on an initial approximation, and are quite sensitive to that initial approximation. In addition, the assumption that the internal camera parameters are fixed is not always a good assumption.\nAnother set of methods does not require determining the values of the various parameters in equation (1). Instead, reconstruction is performed from an examination of various features of the scene which are present in images of the scene that are recorded from a plurality of diverse locations. One such method, described in U.S. Pat. No. 5,598,515, entitled \"System And Method For Reconstructing Surface Elements Of Solid Objects In A Three-Dimensional Scene From A Plurality Of Two Dimensional Images Of The Scene,\" issued Jan. 28, 1997, in the name of Amnon Shashua, describes a system and method for scene reconstruction from two views. In the system and method described in that patent, the epipole for at least one of the images is determined, which can lead to error in the reconstruction. In several other methods, described in, for example, A. Shashua, \"Algebraic functions for recognition,\" IEEE Transactions on Pattern Analysis and Machine Intelligence, Vol. 17, No. 8, pp. 779-789 (1995), and R. Hartley, \"Lines and points in three views-a unified approach,\" Proceedings of the 1994 Image Understanding Workshop, pp. 1009-1016, do not require the direct determination of the epipoles, but those methods require use of three views. All of these methods require that corresponding points and/or lines, that is, points and/or lines in the views which are projections of the same points and/or lines in the scene, be located in all of the three views. In some applications, locating corresponding points and/or lines in the three views can be difficult or impossible."}
{"text": "The present invention relates to an infant care apparatus and, more particularly, to an infant incubator having a uniform and symmetrical flow pattern to provide an even distribution of heat within the infant compartment.\nThere are, of course, many infant incubators available today for the caring of infants and to provide heat and, generally, humidity to that infant to aid in its wellbeing. In general, such incubators comprise an infant platform that is affixed to a base so as to locate the infant platform at a convenient height for the caregiver to view and attend to the infant. Within the infant platform, there is a flat, planar surface that underlies and supports the infant. A plurality of transparent sides generally extend upwardly from the infant platform to surround the infant and create the isolated infant compartment within which the infant is positioned. As indicated, the infant compartment allows the creation of a controlled environment and provides heat and humidity to the infant to aid in its growth and development.\nIn most incubators, the infant is heated by means of a heated flow of air that is introduced into the infant compartment to warm the infant and there is a recirculation of air from the infant compartment to be reintroduced to the heating system to be re-circulated back into the infant compartment. Thus, there is a heating and air moving compartment, normally positioned within the infant platform and located underneath the flat, planar surface on which the infant lies so as to be close to the infant compartment in order to conserve heat to the infant. That heating and air moving compartment contains the necessary components to heat and force the air into the infant compartment and conventionally comprise an electric heater and a motor powered fan to provide that heated air to the infant compartment.\nThe actual introduction of the heated air and the removal of air from the infant compartment also varies among incubators, however, it is preferable that the air flow within the infant compartment be uniform so as to create a uniform, even heat distribution within that compartment. As used hereinafter, since the base of the infant compartment is generally rectangular, there is a long dimension that will be referred to as the longitudinal side of the infant compartment or incubator and a shorter dimension that will be referred to as the end of the infant compartment or incubator. In use, the infant is normally placed parallel to the longitudinal sides of the infant compartment and there are access doors along those sides or side, so that access can be had to the infant by the caregiver.\nIn various incubators there are certainly differing systems to introduce the heated air into the infant compartment. For example, in U.S. Pat. No. 5,730,355 of Lessard et al, there is a flow pattern that introduces the heated air into an infant compartment through openings along the elongated sides of the infant compartment and to cause the air to exit from the infant compartment along only one of the ends of the infant compartment. Accordingly, the air flow is not uniform as all of the air introduced is removed along one of the ends and would tend to create a cooler area at the upper end of the infant compartment where there is no opening for the air to exit the infant compartment. The non-uniformity can also create a cooler area on the inside walls of the incubator.\nFurther, in Beer et al, U.S. Pat. No. 4,617,912, there is an air flow pattern wherein the heated air enters the infant compartment along both of the sides of the incubator and exits the infant compartment along the ends of the infant compartment. While more uniform, it is noted that the air directly enters the infant compartment at the bottom of the compartment and does not pass through a double walled door to enter the infant compartment at the top thereof. Thus, the inflow of warm air is subject to partial blockages or diversions that may be caused by a blanket or other material within the infant compartment that can upset the uniformity of the warm air flow at the sides of the incubator where such items may be present.\nAs an example of means to avoid partial blockages of side openings of the warm air into an infant compartment, double walls are used in some infant incubators, such as in Koch et al U.S. Pat. No. 4,936,824, where the inner walls of the infant compartment are actually heated by the flow of the heated air to raise the temperature of those walls in order to reduce the radiant energy loss of the infant to such walls. When such double walls are used, therefore, it is desirable that the wall temperature be fairly uniform across the inner surface of the walls. Again, therefore, a uniform, symmetrical flow pattern would be advantageous. As is noted in Koch et al, while blockage of the inlet flow paths are eliminated, there is a distribution pattern of the heated air that follows the inlet passage within the hood from front to back of the infant compartment where the air is substantially contained within the walls of the hood. Basically, in the Koch et al construction, there is an inlet along one longitudinal side and one outlet along the opposite longitudinal side with the heated air contained generally within a passageway formed between a double walled hood.\nAccordingly, it would be desirable to provide a heating system that introduces the heated air into the infant compartment in a uniform pattern so that the areas within the infant compartment are generally at the same temperature and with the use of air passages through double walls of the infant compartment. As such, the areas of the double walls are at a more uniform temperature to lower the heat loss of the infant. In addition, it would be further advantageous to be able to provide a uniform flow pattern with a minimum of components in convective heating system by means of a heating and air moving compartment that is efficient and yet which provides the heating and air movement along the paths desired by utilizing large, molded components.\nAccordingly, the present invention relates to an infant incubator that has a base and a frame extending upwardly from the base and which has an infant platform extending outwardly, preferably in cantilever manner, from the frame. The infant platform has a flat, planar, generally rectangular surface to provide a support for positioning the infant and has longitudinal sides and ends. An infant compartment is formed about the infant resting on the flat, planar surface by means of peripheral walls extending upwardly from said longitudinal sides and the ends and a cover that encloses and retains the infant within a controlled environment.\nA convective heating system is used to force a flow of heated air into the infant compartment and to return air from that infant compartment. In the preferred embodiment, the components of the convective heating system are located in a heating and air moving compartment that is located within the infant platform just beneath the flat, planar surface on which the infant rests to minimize the loss of heat between the convective heating system and the infant compartment. Those components of the convective heating system include a fan, a motor to power the fan, a heater, and the various ducting and passageways to direct the air into and out of the infant compartment.\nIn that heating and air moving compartment, the convective heating system is readily constructed and assembled by a minimum, of large molded plastic components including a one-piece molded chassis and a molded upper pan that fits over the chassis and forms therebetween a lower chamber and an upper chamber.\nIn the preferred system, elongated inlets are formed along the opposite ends of the infant compartment and allow air to pass from the infant compartment through those elongated inlets and which then communicate with the lower chamber. In similar fashion, the upper chamber communicates with the infant compartment through elongated outlets that are formed along the longitudinal opposite sides of the infant compartment. A fan, powered by a motor, is located intermediate the upper and the lower chambers and serves to force the air from the lower chamber to the upper chamber and thus, to move the air from the infant compartment via the elongated inlets to return to the infant compartment t through the elongated outlets. In the preferred embodiment, a heater is located in the lower chamber to heat the air passing through that chamber.\nIn the system, therefore, the air thus moves from the infant compartment along the opposite ends and returns to the infant compartment along the opposite longitudinal sides of the infant compartment. Thus, with the use of only two molded plastic components, the system efficiently heats and delivers the air as desired with a minimum of components and thus is readily assembled.\nBy means of the present invention, the heated air from the convective heating system is introduced into the infant compartment along both of the longitudinal sides of the infant compartment and the return air leaves the infant compartment along both of the ends of the infant compartment via openings along those respective longitudinal sides and ends of the infant platform.\nIn the present invention, the side longitudinal walls of the infant incubator are double walls, that is, they comprise an inner and an outer wall spaced apart so as to form a passageway between the inner and outer walls to allow the heated air to pass through the longitudinal walls such that the heated air is introduced into the infant compartment at or near the upper regions of the infant compartment. In such manner, the side longitudinal walls are heated much in the same manner as in the Koch et al U.S. Pat. No. 4,936,824 to resist radiant losses from the infant through those walls.\nAccording, the flow of heated and return air is basically symmetrical and uniform through the infant compartment. The heated air enters the infant compartment along both of the longitudinal sides of the flat, planar surface and leaves the infant compartment through both of the ends of the infant compartment and therefore the flow pattern, as well as the heat distribution throughout the infant compartment, is more uniform.\nThese and other features and advantages of the present invention will become more readily apparent during the following detailed description taken in conjunction with the drawings herein."}
{"text": "The technical field of this invention is methods for screening for transdominant effector peptides and RNA molecules selected inside living cells from randomized pools.\nSignaling pathways in cells often begin with an effector stimulus that leads to a phenotypically describable change in cellular physiology. Despite the key role intracellular signaling pathways play in pathogenesis, in most cases, little is understood about a signaling pathway other than the initial stimulus and the ultimate cellular response.\nHistorically, signal transduction has been analyzed by biochemistry or genetics. The biochemical approach dissects a pathway in a xe2x80x9cstepping-stonexe2x80x9d fashion: find a molecule that acts at, or is involved in, one end of the pathway, isolate assayable quantities and then try to determine the next molecule in the pathway, either upstream or downstream of the isolated one. The genetic approach is classically a xe2x80x9cshot in the darkxe2x80x9d: induce or derive mutants in a signaling pathway and map the locus by genetic crosses or complement the mutation with a cDNA library. Limitations of biochemical approaches include a reliance on a significant amount of pre-existing knowledge about the constituents under study and the need to carry such studies out in vitro, post-mortem. Limitations of purely genetic approaches include the need to first derive and then characterize before proceeding with identifying and cloning the gene.\nScreening molecular libraries of chemical compounds for drugs that regulate signaling systems has led to important discoveries of great clinical significance. Cyclosporin A (CsA) and FK506, for examples, were selected in standard pharmaceutical screens for inhibition of T-cell activation. It is noteworthy that while these two drugs bind completely different cellular proteinsxe2x80x94cyclophilin and FK506 binding protein (FKBP), respectively, the effect of either drug is virtually the samexe2x80x94profound and specific suppression of T-cell activation, phenotypically observable in T cells as inhibition of mRNA production dependent on transcription factors such as NF-AT and NF-xcexaB. Libraries of small peptides have also been successfully screened in assays for bioactivity. The literature is replete with examples of small peptides capable of modulating a wide variety of signaling pathways. For example, a peptide derived from the HIV- 1 envelope protein has been shown to block the action of cellular calmodulin.\nA major limitation of conventional in vitro screens is delivery. While only minute amounts of an agent may be necessary to modulate a particular cellular response, delivering such an amount to the requisite subcellular location necessitates exposing the target cell or system to relatively massive concentrations of the agent. The effect of such concentrations may well mask or preclude the targeted response.\nThe present invention provides methods and compositions to create, effectively introduce into cells and screen compounds that affect a signaling pathway. Little or no knowledge of the pathway is required, other than a presumed signaling event and an observable physiologic change in the target cell. The disclosed methods are conceptually distinct from prior library search methods in that it is an in vivo stratagem for accessing intracellular signaling mechanisms. The invention also provides for the isolation of the constituents of the pathway, the tools to characterize the pathway, and lead compounds for pharmaceutical development.\nRelevant Literature\nMann et al. (1983) Cell 33, 153-159, Pear et al. (1993) Proc. Natl. Acad. Sci. USA 90(18):8392-6 and copending U.S. patent application Ser. No. 08/023,909 describe the BOSC and BING retroviral systems useful as delivery vectors for the disclosed methods.\nScott and Craig (1994) Current Opinion in Biotechnology 5:40-48 review random peptide libraries. Hupp et al. (1995) describe small peptides which activate the latent sequence-specific DNA binding function of p53. Palzkill et al. (1994) report the selection of functional signal cleavage sites from a library of random sequences introduced into TEM-1-lactamase.\nThe invention provides methods and compositions for screening for transdominant intracellularly bioactive agents such as pharmaceuticals. The invention accesses molecules or targets within living cells and provides for the direct selection of those bioactive agents with desired phenotypic effects. The general methods involve steps: expressing a molecular library of randomized nucleic acids as a plurality of isolated corresponding randomized expression products in a plurality of cells, each of the nucleic acids comprising a different nucleotide sequence, screening for a cell of the plurality of cells exhibiting a changed physiology in response to the presence in the cell of a transdominant expression product of the corresponding expression products; and detecting and isolating the cell and/or transdominant expression product,\nIn a particular embodiment, the expressing step comprises translating the nucleic acids and/or corresponding transcripts, and each of the nucleic acids encodes a peptide comprising a different amino acid sequence. The nucleic acids may be joined to sequences encoding polypeptide backbones of artificial design capable of intracellularly presenting randomized peptides as structured domains. The methods may also involve introducing the library into the cells, such as through the use of retroviral vectors, and particularly suitable vectors are also disclosed."}
{"text": "1. Technical Field\nThe present invention relates to a passive intermodulation (PIM) measurement instrument configured to utilize frequency-dependent plugins.\n2. Related Art\nA PIM is an unwanted signal or signals generated by the non-linear mixing of two or more frequencies in a passive device such as a connector or cable. PIM has surfaced as a problem for cellular telephone technologies such as Global System for Mobile Communications (GSM), Advanced Wireless Service (AWS) and Personal Communication Service (PCS) systems. Cable assemblies connecting a base station to an antenna on a tower using these cellular systems typically have multiple connectors that cause PIMs that can interfere with system operation.\nThe PIM signals are created when two signals from the same or different systems mix at a PIM point such as a faulty cable connector. If the generated PIM harmonic frequency components fall within the receive band of a base station, it can effectively block a channel and make the base station receiver think that a carrier is present when one is not. PIMs can, thus, occur when two base stations operating at different frequencies, such as an AWS device and a PCS device, are in close proximity.\nThe PIMs can be reduced or eliminated by replacing faulty cables or connectors. Test systems are thus utilized to detect the PIMs enabling a technician to locate the faulty cable or connector. The test system to measure the PIMs, thus, creates signals at two different frequencies, amplifies them, and provides them through cables connecting a base station to antennas on a tower for the base stations. A return signal carrying the PIMs is filtered to select a desired test frequency harmonic where PIMs can be detected and the PIM measurement is provided to an operator.\nPIM testers to date have used CW signals for the two frequencies used to create the PIM. This is due to the unknown nature of where physically the PIM is located in the transmission path. The PIM is monitored by one technician while the other technician climbs the tower and physically moves the connector joints to see if the PIM changes. Other techniques plot a time graph of the PIM so a single technician can correlate his movement up the tower with results on a graph provided on a plotter below the tower.\nFIG. 1 shows a block diagram of components of a prior art test system setup for measuring a PIM. The test system utilizes two signal sources 100 and 102 producing CW signals, with a first signal source 100 producing a signal at frequency F1 and the second signal source 102 producing a signal at frequency F2. When these multiple signals are allowed to share the same signal path in a nonlinear transmission medium, the unwanted signals can occur. The combined 3rd order response is particularly troublesome as it produces an unwanted signal at 2F1-F2 that can pass from one system transmitter into another system's receiver.\nThe signal at frequency F1 is provided from source 100 to a high power amplifier (HPA) 104. The signal at frequency F2 is provided from source 102 to a high power amplifier 106. Both the high power amplifiers 104 and 106 are shown as 50 W amplifiers, and receive a DC power supply input shown ranging from 100 to 125 Watts to produce a 50 Watt signal output.\nThe output of each of the amplifiers 104 and 106 is provided through respective isolators 108 and 110 to the input of a hybrid combiner 112. The hybrid combiner 112 assures the two carrier signals F1 and F2 are isolated from each other. If they are allowed to combine without isolation, intermodulations would appear due to power output stage nonlinearities. The isolators 108 and 110 are inserted after the power amplifiers 104 and 106 to give additional isolation from any return signal from the hybrid combiner 112. The intermodulations are the same frequency as the PIM (2F1-F2), so isolation using both the hybrid combiner 112 and the isolators 108 and 110 is critical.\nThe outputs of the hybrid combiner 112 are provided to a commercial duplexer 114. The commercial duplexer 114 includes a transmit filter 116 and a receive filter 118. Signals F1 and F2 are provided to a first terminal via the transmit filter 116, the first terminal is connected to the front panel 120 of the instrument using a Low PIM Cable 122. Although the low PIM cable 122 is designed to minimize PIM, the low PIM cable can introduce additional PIM signals into the system. The PIM signal is provided to a second terminal via the receive filter 118. The PIM signal can be provided to a digital receiver or spectrum analyzer for measurement, such as a tuned receiver 124 which is connected to the second terminal via a FPIM filter 126 and Pre Amp 128.\nThe power needed to create the PIM is a standardized 20 W per carrier. Overall for the PIM test circuit of FIG. 1, the DC power supplied to the amplifiers 4 and 14 needs to be 50 to 60% higher to create the 20 W output due to DC and RF inefficiencies. This translates to a continuous DC power consumption of 200 to 250 Watts. The loss of power combiner 20 is 3 dB, so 25 W carriers (F1 and F2) can emerge from the combiner while other 25 W carriers that are not needed are dissipated in an internal load 130. The 50 Watt power output of the two amplifiers 104 and 106 is further reduced a total of at least 1 dB above the theoretical 3 dB loss through the hybrid combiner 112 due to the losses through cabling. Further losses in the isolators 108 and 110 and commercial duplexer 114 reduce total power so that 20 Watt carriers F1 and F2 are produced from the output of commercial duplexer. PIMs introduced by the low PIM cable 122 or other sources receiving the signals F1 and F2 will generate a return PIM signal that is provided back through commercial duplexer 114 and directed to the tuned receiver 124 for processing."}
{"text": "With increases in aging population, the development of pharmaceuticals effective for the treatment of senile dementia has been demanded strongly in recent years. Alzheimer's disease, which is a typical disease of senile dementia, is a neurodegenerative disease characterized by brain shrinkage, deposits of senile plaques, and neurofibrillary changes. This disease develops when insolubilized molecules attributed to a change of forms of amyloid β peptides and following fibrillogenesis are deposited on neuronal cells such that neuronal cell death is induced by their toxicity.\nAmyloid β peptides (Aβ) are degradation products generated by cleavage with β secretase, etc., from neuronal amyloid precursor proteins, and these degradation products are found in 2 forms, Aβ1-40 and Aβ1-42. It has been reported that Aβ1-42 is more likely to aggregate and more correlated to disease and neurotoxicity.\nIf a change of forms of amyloid β peptides and fibrillogenesis thereof can be inhibited, the development of Alzheimer's disease can be prevented.\nThe pamphlet of WO 2005/105998 describes that a single-chain antibody having an activity of specifically binding to Aβ1-42 such that the fibrillogenesis reaction thereof is inhibited is useful as a preventive or therapeutic agent for Alzheimer's disease.\nHowever, antibodies are high-molecular-weight proteins and are therefore expensive. They also have problems such as complicated production and purification steps or a lack of stability.\nDrugs having an effect of suppressing amyloid β peptide production, for example, a rhodanine derivative (JP Patent Publication (Kokai) No. 06-192091A (1994)), a benzimidazole derivative (U.S. Pat. No. 5,552,426), a vinpocetine derivative (pamphlet of WO 96/25161), and an aromatic amide derivative (pamphlet of WO 2004/014843), have been known as low-molecular-weight preventive or therapeutic agents for Alzheimer's disease."}
{"text": "It is necessary in hospitals and other medical environments to organize medical implants and surgical instrumentation in surgical cases. The implants and instruments are typically held in various brackets and trays within the case, so as to be presented in a convenient, organized fashion upon their use in a medical procedure, such as a surgical procedure. These surgical cases are typically provided with a lid to protect the contents.\nThe demands on medical personnel involved in surgical procedures are often very great. For example, many procedures may be performed in the context of trauma treatment, requiring prompt treatment of an injured patient. Even in the context of the scheduled treatment of a patient, unexpected developments during the course of the procedure can create great urgency. Accordingly, medical personnel are often operating under substantial time constraints and a great deal of resultant psychological pressure, based on the stakes involved in the proper performance of their duties (the health, or possibly even the life, of the patient). It is not an exaggeration to state that the delay of even a few seconds can cost a life in emergency treatment. Therefore, it is critical that the apparatus and methods used by such personnel are designed to work as efficiently and safely as possible, in order to allow the personnel to perform their work quickly, as well as to accommodate minor mistakes and mishaps which may occur in the course of such high-pressure work.\nIn the context of surgical cases, a case must offer a convenient hand-hold so that it can be quickly grasped and lifted. The hand-hold may be a handle protruding from the case body, an opening in the case walls for the fingertips to be inserted, a protrusion on the case walls for fingertip gripping, or in some cases may simply be the flat outer side walls of the case itself.\nIdeally, the hand-hold will allow the case to be grasped and lifted by one hand. This will allow medical personnel who are already using one hand to carry something else, to grasp and lift the case. Also, a case which is designed to be carried by one hand should accommodate the situation where it is grasped and lifted by both hands, and then one hand loses its grip. It is also preferable that at least when not in use, the hand-hold will not protrude too far from the walls of the case.\nIt is also common for surgical trays to be presented in situations where physical access to the tray is limited. For example, trays may often be stored in such a way that lateral access to the ends or side of the tray is not available. Furthermore, surgical trays may be disposed in use in sterilization enclosures which also limit lateral access to the tray. In such situations, grasping and lifting of prior art trays, in which the handles are projections on the end walls of the base or around the edges of the top of the lid, will be difficult if not impossible. Accordingly, there is a need in the art for a system which provides a handle system allowing convenient grasping and lifting of the surgical tray when the tray can only be accessed from above.\nPrior art hand-hold systems for surgical cases do not typically meet all of these design criteria. For example, U.S. Pat. No. 4,728,504 to Nichols, entitled \"Stackable Medical Instrument Sterilizer Container\" discloses a container 10 including a housing 12 having handles 28 provided on opposite ends of the container to facilitate handling of the container. The handles 28 are arched protrusions which allow only a fingertip grip, which is weaker than a full hand grip. Also, these handles 28 do not permit the case to be securely carried by one hand. U.S. Pat. No. 5,415,846 to Berry, Jr., entitled \"Plastic Autoclave Tray and Lid Combination\" discloses an autoclave tray 21 having an outwardly extending flange 25 and is typical of prior art systems which have a simple rim around the edge of the case. The rim can be used for a fingertip hold, and has the same problems as discussed above with reference to the Nichols '504 patent. Handle systems such as those of the Berry and Nichols patents do not allow the tray to be lifted from above when the tray system is in a confined space such that only the top surface of system is accessible. Specifically, if those trays were in an enclosure open at the top and closely matching the outer periphery of the tray, they could not be grasped and lifted, as it would not be possible for the user's hands to grip the handles.\nThe lid of a surgical case must also be designed with both safety and ease of use in mind. A lid which is unduly cumbersome to remove will cause an unacceptable delay in accessing the products inside the case. On the other hand, if the lid comes off too easily, it may come off accidentally or prematurely, for example during transport. A loose lid that accidentally becomes separated from the case may allow the case contents to become disarranged or result in the spillage of the case contents, delaying the surgery and possibly resulting in damage to the case contents. Also, a loose lid may in itself present a hazard if it falls to the floor in a busy area.\nUnderlying all of the constraints discussed above are fundamental ergonomic design concepts which apply to any product which people use. A product should be a simple as possible while accomplishing the required function, with a minimum of moving parts, and its method of use should be simple and easily understood from the product itself. Furthermore, the product should be simple and inexpensive to manufacture. Thus, it is preferable that the hand-hold and lid-securing functions be combined, such that they are performed at least in part by the same structural elements of the case system. The simplicity of the mechanical design is particularly important in the surgical tray field, as complex designs having many mating parts will be more difficult to properly sterilize.\nSystems combining the hand-hold and lid-securing functions are known in the art. For example, U.S. Pat. No. 4,915,913 to Williams et al., entitled \"Medical Sterilizer Device with Improved Latch Mechanism\" discloses a medical sterilization case 4 and a cover C which is latched to the body of case by a pair of latch mechanisms 20. Each latch mechanism includes a handle 22 and a complicated latch mechanism 20 which includes, among other components, a hinge plate 106, an elongated retention member 110, a second retention member 116, a latch plate 130, and coil springs 142. These components are arranged so as to provide a pivoting, double-action latching system for the cover C. The Williams system is unduly complex for the simple functions being performed. Also the handle 22 is provided as a means to accomplish the unlatching of the latch mechanism, and is not well adapted to be used to pick up the case. If used to pick up the case, the force applied to the handle 22 would tend to result in the inadvertent unlatching of the latch mechanism, undesirably allowing the cover to come free prematurely.\nThe Berry '846 patent, discussed above, also discloses a relatively complex latching mechanism. In Berry, a frame 50, a spring clip 51, a slide 52, collar 53, and cooperating screw stud 53a are assembled to form a slide latch assembly which allows the lid 22 to be releasably secured to the tray 21. Similarly, U.S. Pat. No. 5,384,103 to Miller, entitled \"Instrument Tray,\" also discloses a complex locking mechanism 16 which is disposed on the outer surface of each of the tray end walls (see, for example, FIG. 4 of Miller and associated text at 6:54-7:42). Also, the handles of the Miller tray may be accessed through openings in the cover, but the system is not designed to take advantage of any mechanical association of the handles with the lid to contribute to securement of the lid. Although Miller is directed to the problem of providing a releasable locking system for a tray lid, as in the Williams and Berry patents, the locking mechanism 16 provided by Miller is an inordinately complex design solution for the problem addressed.\nAccordingly, there is a need in the art for a combined surgical case handling and lid securing system having improved ergonomics and allowing a firm, secure full-hand grasp, and which permits one-handed grasping, lifting, and carrying of the case as well as positively engaging the lid to the case until its removal is desired. The handle of the lid is accessible even when the system is approachable only from the top. The lid should remain positively secured to the case even when the case is held sideways, for example when one hand is used to grasp, lift and carry the case. The hand-holding and lid-securing functions should be combined in a single system, such that parts serving one function also contribute to the other function. There is also a need in the art for a surgical case lid-securing system meeting these specifications that is simple and inexpensive to manufacture and is simple to use."}
{"text": "The present disclosure relates generally to information handling systems, and more particularly to programming fuses in chipsets used in information handling systems.\nAs the value and use of information continues to increase, individuals and businesses seek additional ways to process and store information. One option available to users is information handling systems. An information handling system generally processes, compiles, stores, and/or communicates information or data for business, personal, or other purposes thereby allowing users to take advantage of the value of the information. Because technology and information handling needs and requirements vary between different users or applications, information handling systems may also vary regarding what information is handled, how the information is handled, how much information is processed, stored, or communicated, and how quickly and efficiently the information may be processed, stored, or communicated. The variations in information handling systems allow for information handling systems to be general or configured for a specific user or specific use such as financial transaction processing, airline reservations, enterprise data storage, or global communications. In addition, information handling systems may include a variety of hardware and software components that may be configured to process, store, and communicate information and may include one or more computer systems, data storage systems, and networking systems.\nMany information handling systems include chipsets such as, for example, the Platform Controller Hub (PCH) available from INTEL® Corporation of Santa Clara, Calif., United States. Such chipsets may include One-Time-Programmable (OTP) Non-Volatile Memory (NVM) and/or other memory subsystems that utilize fuses such as In Field Programmable (IFP) fuses that may be burned during manufacture to provide security information in the chipset (e.g., information associated with a public key such as a hash of a master public key) and/or enable particular features in the system. For example, security information provided in the chipset in such a manner and associated with a public key may be utilized to verify software or firmware (e.g., a Basic Input/Output System (BIOS)) that has been signed with an associated private key. Similarly, features enabled in the information handling system in such a manner may include security features such as, for example, BOOT GUARD available in systems provided by DELL® Inc. of Round Rock, Tex., United States, and PLATFORM TRUST TECHNOLOGY available from INTEL® corporation of Santa Clara, Calif., United States. Providing chipsets with information and enabling system features in such manner provides a root of trust for that information and those features.\nConventional systems provide two methods for programming IFP fuses during manufacture: Automatic End Of Manufacturing (Auto EOM) and Host Embedded Controller Interface (HECI) EOM. Auto EOM provides for the burning of fuses in the chipset automatically upon the initial boot of the system and before the BIOS begins executing. As such, Auto EOM cannot be controlled by the BIOS or any software tool or script. HECI EOM provides for the burning of the fuses in response to a command provided post-boot (e.g., using an SPSManuf Universally Extensible Firmware Interface (UEFI) shell tool or executable version that can be run in an operating system.) Neither of Auto EOM or HECI EOM differentiates between the chipset (e.g., the PCH) or the Central Processing Unit (CPU) installed in the system, which can raise some issues. For example, each of PCH and CPU may be provided in different versions (e.g., a production version (also referred to as a Qualification Sample (QS)) or a pre-production version (also referred to as an Engineering Sample (ES)). The IFP fuse burn operations may differ based on the combination of the PCH and the CPU utilized in a system, and conventional BIOS storage only stores a single set of IFP fuse burn instructions. As such, a manufacturer must determine and provide the appropriate IFP fuse burn instructions in the BIOS storage for use in the IFP fuse burning process, and risks rendering the system unusable (e.g., “bricked”) if the wrong IFP fuse burn profile is utilized.\nAccordingly, it would be desirable to provide an improved chipset fuse programming system."}
{"text": "Wavelength division multiplexing (WDM) is one technology that is envisioned to increase bandwidth capability and enable bi-directional communications in optical networks. In WDM networks, multiple data signals can be transmitted simultaneously between network elements (NEs) using a single fiber. Specifically, the individual signals may be assigned different transmission wavelengths so that they do not interfere or collide with each other. The path that the signal takes through the network is referred to as the lightpath. One type of WDM network, a wavelength switched optical network (WSON), seeks to switch the optical signals with fewer optical-electrical-optical (OEO) conversions along the lightpath, e.g. at the individual NEs, than existing optical networks. One of the challenges in implementing WSONs is the determination of the routing and wavelength assignment (RWA) for the various signals that are being transported through the network at any given time. To implement RWA, various NE related information can be forwarded from a Path Computation Client (PCC), such as a NE, and received and processed at a Path Computation Element (PCE)."}
{"text": "In certain applications such as vehicle exhaust systems it is desired to provide a pipe joint with high pull-apart strength and with a good fluid seal between the pipes. It has become increasingly important, especially in connection with vehicle exhaust systems, to achieve greater reliability and uniformity in providing pipe couplings with a good fluid seal.\nIn the prior art, a pipe coupling which provides high strength and a good fluid seal has been provided by a pipe lap joint having a collapsible sealing zone and a band clamp as disclosed and claimed in Cassel U.S. Pat. No. 4,629,226 granted Dec. 16, 1986. In the lap joint of that patent, a collapsible sealing zone is provided in the outer pipe by at least one pair of intersecting end-to-end slots which afford relief for circumferential contraction of the sealing zone within the overlap region of the pipe ends. A clamping band is disposed around the outside pipe and covers the inboard slot. The clamping band is tightened around the outside pipe to clamp the pipes together and to collapse the sealing zone of the outer pipe into close fitting engagement with the inner pipe and thereby cause sealing engagement of the outer pipe with the inner pipe and with the band. The intersecting slots comprise an outboard slot and an inboard slot. The outboard slot has sidewalls extending inwardly from the end of the outside pipe and terminating in a transverse inner end wall within the overlap region of the pipes. The inboard slot is circumferentially offset from the outboard slot with sidewalls extending axially inwardly from an outer end wall. The adjacent sidewalls of the slots are in substantial alignment with each other and the end walls are in substantial alignment with each other. When the overlap region of the outer pipe is collapsed into close fitting engagement with the inner pipe, the end walls of the slots are moved into engagement with each other and thereby isolate the slots from each other and provide a seal.\nThe lap joint of the aforementioned U.S. Pat. No. 4,629,226 utilizes a pair of slots which intersect each other at adjacent corners with the outer pipe in an uncollapsed state. The inner pipe and outer pipe walls are telescoped together in a loose fit to form the lap joint. The manufacturing tolerances of the pipe diameters result in a large percentage of the mated pipes having a sufficiently loose fit that the outer pipe is collapsed sufficiently by the band clamp that the end walls of the slots are moved into engagement with each other to isolate the inboard and outboard slots. However, there is some small percentage of mating pipes in which the initial fit is so close that tightening of the band clamp causes very little collapse of the outer pipe before it is in close fitting engagement with the inner pipe and there is insufficient displacement to cause the edge-to-edge engagement of the end walls of the slots which is necessary to provide isolation between the two slots. This may lead to undesirable leakage in some of the pipe joints.\nIn the earlier prior art, it was known to use a collapsible sealing zone at the free end of the outer pipe which is collapsed into sealing engagement with the inner pipe by a clamping band. A pipe lap joint of this type is disclosed in Wagner et al. U.S. Pat. No. 4,113,289 granted Sep. 12, 1978. The collapsible sealing zone in the Wagner et al. patent comprises several sets of slots which are disposed circumferentially around the end of the outer pipe. In each set of slots there is an open slot which extends inwardly from the end of the pipe and an adjacent closed slot, i.e. one which does not extend to the end of the pipe. The difficulty with this sealing zone is that the slot structure exhibits a high resistance to collapsing and accommodates a relatively small amount of reduction in the circumference of the pipe at each set of slots.\nA lap joint of the type using a collapsible sealing zone is also disclosed in Cassel U.S. Pat. No. 4,056,273. In the joint of this patent, a collapsible sealing ring forms an extension of the outer pipe but is separate from it. The sealing ring is split so as to form a tongue on each free end with the tongues having complementary ramp surfaces in engagement with each other and which slide relative to each other when the ring is contracted or collapsed. This sealing ring provides a good seal but the structure is not adapted to a sealing ring which is integral with the pipe end.\nA general object of this invention is to overcome certain disadvantages of the prior art and to provide a pipe lap joint with an improved sealing zone."}
{"text": "Some 1,2,4-triazolo[1,5-c]pyrimidines are known to the art. Certain 1,2,4-triazolo[1,5-c]pyrimidines are disclosed as being bronchodilators in the patents discussed below, the compounds being referred to therein as triazolo[2,3-c]pyrimidines:\nUnited Kingdom Patent No. 859,287 discloses 2-amino-1,2,4-triazolo[1,5-c]pyrimidines which are substituted on the pyrimidine ring at the 5, 7 and 8 positions by certain combinations of substituents selected from hydrogen, alkyl, halogen-substituted alkyl, hydroxy-substituted alkyl, alkoxy-substituted alkyl, alkenyl, cycloalkyl, amino, alkylamino, dialkylamino, phenyl, alkylthio, alkoxy and halogen substituents. United Kingdom Patent No. 898,407 discloses processes for preparing certain of these compounds by subjecting the corresponding 1,2,4-triazolo[4,3-c]pyrimidines to an acid treatment, to an alkaline treatment, or to a heat treatment.\nUnited Kingdom Patent No. 873,223 discloses 2-amino or 2-acetamido-1,2,4-triazolo[1,5-c]pyrimidines which are substituted on the pyrimidine ring at the 5, 7 and 8 positions by certain combinations of substituents selected from hydrogen, alkyl, halogen-substituted alkyl, alkoxy-substituted alkyl, alkenyl, cycloalkyl, alkylthio and halogen substituents.\nUnited Kingdom Patent No. 897,870 discloses 2-alkylamino-1,2,4-triazolo[1,5-c]pyrimidines, 2-dialkylamino-1,2,4-triazolo[1,5-c]pyrimidines, and 1,2,4-triazolo[1,5-c]pyrimidines containing a piperidino or morpholino substituent bonded at the 2-position through the nitrogen atom, which compounds are substituted on the pyrimidine ring at the 5, 7 and 8 positions by certain combinations of substituents selected from hydrogen, alkyl, halogen-substituted alkyl, hydroxy-substituted alkyl, alkenyl and halogen substituents.\nThe following related articles discloses the synthesis of certain 1,2,4-triazolo[1,5-c]pyrimidines as potential bronchodilators:\nG. W. Miller et al., J. Chem. Soc., 1963, 5642, discloses 2-amino- or 2-acetamido-1,2,4-triazolo[1,5-c]pyrimidines (referred to therein as triazolo[2,3-c]pyrimidines) which are substituted on the pyrimidine ring by, for example, hydrogen and alkyl substituents. Certain of these compounds are said to be bronchodilators.\nG. W. Miller et al., J. Chem. Soc., 1963, 3357, discloses 1,2,4-triazolo[1,5-c]pyrimidines (referred to therein as triazolo[2,3-c]pyrimidines) which are substituted at the 2-position by hydroxy, halogen, alkoxy, amino or substituted amino substituents and on the pyrimidine ring by alkyl substituents, or alkyl and halogen-substituted alkyl substituents.\nW. Broadbent et al., J. Chem. Soc., 1963, 3369, discloses 1,2,4-triazolo[1,5-c]pyrimidines (referred to therein as triazolo[2,3-c]pyrimidines) which are substituted at the 2-position by a mercapto, alkylthio, alkylsulphonyl, or dialkylamino substituent, and on the pyrimidine ring by alkyl substituents or alkyl and halogen-substituted alkyl substituents.\nStill other 1,2,4-triazolo[1,5-c]pyrimidines are disclosed in the following articles and patents.\nTemple et al., J. Org. Chem., 1963, 33, 530, discloses the compound 8-amino-7-chloro-s-triazolo[1,5-c]pyrimidine-2(3H)-one.\nD. J. Brown et al., Aust. J. Chem., 1978, 31, 2505, discloses 1,2,4-triazolo[1,5-c]pyrimidines which are substituted at the 2-position by hydrogen or an alkyl substituent, and on the pyrimidine ring by hydrogen and/or alkyl substituents.\nD. J. Brown et al., Aust. J. Chem., 1979, 32, 1585, discloses 1,2,4-triazolo[1,5-c]pyrimidines which are substituted at the 2-position by hydrogen or an alkyl substituent, and on the pyrimidine ring at the 5-position by a halogen, hydrazino, alkyl or alkylthio substituent, and at the 7-position by an alkyl substituent.\nD. J. Brown et al., Aust. J. Chem., 1980, 33, 1147, discloses pyrimidines which are substituted at the 2-position by hydrogen or an alkyl or phenyl substituent, and on the pyrimidine ring at the 5-position by halogen, dimethylaminomethyleneamino, hydroxyaminomethyleneamino or 5-acetoxyaminomethyleneamino, and at the 7-position by hydrogen or an alkyl substituent.\nU.S. Pat. No. 4,269,980 discloses 5-, 7- and 8-(optionally substituted phenyl)-1,2,4-triazolo[1,5-c]pyrimidines. These compounds may be substituted at the 2-position by hydrogen or an alkyl substituent and are anxiolytic agents."}
{"text": "The global swine industry has seen increased feeding of by-products (dried distillers grains with solubles (DDGS), wheat midds, etc.) initially from 0-10% to the current extremes of 30-60%. These diet cost savings have been a great opportunity for industry to save on feed input costs, but come with a set of challenges as well. The fermentation process to extract ethanol from corn removes almost all of the starch, leaving the resulting DDGS feed by-product containing approximately 40% fiber. This higher fiber content relative to corn results in reduced dry matter digestibility and approximately 10 percentage units less digestibility of most amino acids in DDGS compared to corn (Stein and Shurson, 2009).\nConsequently the inclusion of DDGS in livestock diets can have negative impacts on animal growth performance and carcass characteristics. In addition to the negative effects on animal growth and carcass quality, alterations in nutrient digestibility as a result of adding DDGS with a high fiber content have implications for swine manure handling, storage, and decomposition. The commercial swine industry has indicated that manure holding capacity is less in anaerobic deep-pit swine manure storage units, and that the manure from pigs fed high level of DDGS has more solids accumulation, as well as ammonia, methane, and hydrogen sulfide gas emissions.\nIn view of the foregoing, it would be desirable to provide Bacillus strains producing enzymes that provide benefits to animals and methods of using these strains."}
{"text": "The present systems and methods relate to processing analyte sensor data from a continuous analyte sensor. More particularly, the present systems and methods relate to providing accurate predictive alarms to a user."}
{"text": "The subject invention relates to a device for use by orthopedic physicians or surgeons in connection with patient bone fracture and ligament management. More particularly, the invention relates to a novel orthopedic cast spacing and spreading device operable to relieve pressure in connection with the use of an immobilizing cast structure.\nWhen a patient fractures an arm or a leg, or damages a ligament to the point that casting is indicated, it is important to insure relative immobility of the member. Accordingly, cast materials are designed to be rigid and do not permit deformation or play in the damaged member.\nIf the soft body tissue of an arm or leg which has been traumatized during an injury, severe enough to fracture a bone and/or damage a joint, swells from the injury within a confined space created by a cast, post casting pressure can cause skin erosion or sores. If such pressure is permitted to persist in muscle compartments, it is possible to actually cause death of the muscle and damage to arteries and/or nerves. This is commonly referred to as compartment syndrome and most references and literature refer to it as \"Volkmann\" disease and contractures. Contractures can result in a significant loss of function of an extremity as well as shortening and cosmetic deformities.\nIn the past, if a surgeon or physician suspected that an extremity was swelling within a cast, the cast was often completely removed and a new cast applied. In many situations, however, it is preferable to simply use a saw to longitudinally cut the cast material and temporarily spread the cast to relieve pressure due to internal swelling. When swelling subsided, the cast could be closed back around the member and rebound with a top layer of plaster cloth.\nA cast spreader device which has previously been known in the art is depicted in a U.S. Pat. No. 3,426,752. The disclosure of this cast spreader instrument is hereby incorporated by reference as though set forth at length for purposes of disclosing the background of the instant invention. Another cast spreader device is disclosed in U.S. Pat. No. 2,757,666. The disclosure of this cast spreader device is also incorporated by reference to further illustrate cast spreading devices existing in the prior art.\nWith conventional plaster cast materials, a cast spreading tool, such as identified above, could be used to separate the cast and the plaster material would remain in a separated or open posture without tending to collapse back to its original configuration. With newer fiber glass and resin materials, however, a physician can spread a cast but as soon as the separating instrument is removed from between the two sides of the cast, the cast to its original configuration tightly around a patient's limb.\nAlthough it has been known, as an expedient, to insert bits of tape or rubber stoppers from syringes into the kerf of a spread apart cast, problems exist with this procedure. More specifically, there is no standard width with such devices and no reliability that pressure on the body member will be relieved. Moreover, such expedient structures have a tendency to fall into a kerf of the cast creating local skin necrosis or further damage.\nThe difficulties suggested in the proceeding are not intended to be exhaustive, but rather are among many which illustrate a need for a device for uniformly providing spacing capability for an orthopedic cast. Other noteworthy problems may also exist, however, those presented above should be sufficient to demonstrate that devices to facilitate orthopedic cast management will admit to worthwhile improvement."}
{"text": "1. Field of the Invention\nThe present invention relates to a computer readable recording medium bearing a printer driver program for editing print job data and a print data processing apparatus in which the printer driver program is installed.\n2. Description of the Background Art\nIn the field of a conventional print data processing system, for example, as disclosed in Japanese Patent laid-open publication Hei 11-1915057, a technique which performs a print operation in accordance with an output paper size designated by changing an output magnification of print job data is proposed, or as described in Japanese Patent laid-open publication 2000-330908, a technique which performs a print operation in accordance with an output paper size edited and designated by a user herself/himself by using an application having a user interface for editing a Markup Language Source is proposed.\nHowever, in a conventional print data processing system, when the print data processing system performs a print operation by using an application which does not have a means for setting and changing the details of a layout such as arrangements of pages, letters, graphics, and the like and does not have a function for editing print data, there is a problem in which a designated print region cannot be printed within an output paper size."}
{"text": "1. Field of the Invention\nThe present invention relates generally to techniques for editing documents or user interfaces. More specifically, the present invention relates to a method and an apparatus that facilitates a grid-based approach to defining table styles.\n2. Related Art\nTables are commonly used to represent information within documents because they can efficiently present information in an intuitive structured format. However, data within tables can sometimes be hard to read. For example, an individual cell within a table can be difficult to distinguish from adjacent cells, especially if the table contains a large number of cells. Furthermore, it is often difficult to trace the column and the row for a given cell back to the corresponding column header and the corresponding row header. Hence, it is often desirable to add attributes, such as cell coloring, to cells within the table which can emphasize and/or distinguish cells from each other, thereby making the contents of the table easier to read.\nTable styles can be defined using several techniques. One such technique involves individually setting attributes for each cell within the table. These attributes can include table-level attributes, row-level attributes, cell-level attributes, and text attributes. Table-level attributes can include table dimensions and stroke/border attributes (e.g. width and color of the table/cell border). Row-level attributes can include row height and keep options. (Keep options specify how paragraphs flow across pages or columns within a document. For example, one keep option keeps a first paragraph and a second paragraph on the same page.) Cell-level attributes can include stroke, fill, diagonal lines, cell insets (i.e. the space between the cell border and text within the cell), vertical justification, and text rotation. Text attributes can include paragraph and character attributes.\nAnother technique for defining table styles involves first selecting a pre-defined pattern, then defining the attributes to use for the pre-defined pattern, and finally applying the pattern of attributes to a target table. For example, Microsoft® Word provides several pre-defined table style patterns which include: whole table, header row, last row, left column, right column, odd row stripes, even row stripes, odd column stripes, even column stripes, top left cell, top right cell, bottom left cell, and bottom right cell. Note that a stripe is a contiguous row or a contiguous column of similar attributes. For example, if the user chooses an odd row stripe pattern, the attributes defined for this pattern is repeated on every cell within an odd-numbered row. Also note that Microsoft® Word enables a user to define a “band” of row/column stripes which allows up to three contiguous rows/columns to contain the same attributes. For example, a user can specify that red is to be applied to alternating rows and that the width of the band should be two rows wide. Hence, the system fills in two rows of red cells, then two rows of white cells (i.e. the default cell background color), then two rows of red cells, etc. in the target table.\nAnother technique involves applying a custom pattern of attributes to a target area within a spreadsheet application. For example, in the Microsoft®Excel® spreadsheet application, a user can define a pattern of color fills which can then be applied to a target area of the spreadsheet. However, there are several limitations to this technique. First, the target area must be an integer multiple of the area (both in x and y directions) containing the pattern to be tiled or repeated. If the target area is not an integer multiple of the area containing the pattern, the pattern is applied to an area within the target area which occupies exactly the same area as the area of the pattern. Second, not all attributes are applied to the target area. For example, whereas color fills and text attributes are applied to the target area, row heights and column widths are not applied to the target area. (Note that in Microsoft® Excel® 2003, the cell formatting can be copied without copying the contents of the area containing the pattern to be tiled.) Third, there is complete lack of control in applying the pattern to the target area. Changing the area containing the pattern to be tiled does not change the target areas where the pattern was applied (or was re-applied)\nUnfortunately, there is no way to generically define a table style which can support all the patterns and attributes for the table or target area. Hence, what is needed is a method and an apparatus for defining and applying table styles to a target table without the problems described above."}
{"text": "1. Technical Field\nThe present disclosure relates to an organic light-emitting diode (OLED) display capable of driving at a low speed and a method of driving the same.\n2. Discussion of the Related Art\nAn active matrix OLED display includes organic light-emitting diodes (OLEDs) capable of emitting light by themselves (i.e., self-emitting), and has advantages of a fast response time, a high emission efficiency, a high luminance, and a wide viewing angle. Each OLED includes an anode electrode, a cathode electrode, and an organic compound layer between the anode electrode and the cathode electrode. The organic compound layer includes a hole injection layer HIL, a hole transport layer HTL, an emission layer EML, an electron transport layer ETL, and an electron injection layer EIL. When a driving voltage is applied to the anode electrode and the cathode electrode, holes passing through the hole transport layer HTL and electrons passing through the electron transport layer ETL move to the emission layer EML and form excitons. As a result, the emission layer EML generates visible light.\nTo reduce power consumption of the OLED display when there is a small change in an input image of the OLED display, a technology for driving the pixels at a low speed has been known. Because a refresh cycle of display data lengthens during a low-speed drive, image quality of the OLED display may be reduced. A main cause of reduction in the image quality is because a voltage level of an emission control signal changes during a data holding period in the low-speed drive.\nMore specifically, the OLED display controls the emission of the pixels through the emission control signal, and the emission control signal is generated by an emission driver. The emission driver may be directly formed in a non-display area (e.g., a bezel area) of a display panel. The emission driver includes a shift register generating an anti-phase emission control signal and an inverter that inverts a phase of the anti-phase emission control signal and generates an emission control signal. Output lines of the shift register for outputting the anti-phase emission control signal are connected to input terminals of the inverter. In this instance, the output lines of the shift register are disposed to intersect clock lines of the inverter to prevent an increase in the size of the bezel area attributable to the emission driver. At least one insulating layer is positioned between the output lines of the shift register and the clock lines of the inverter. Thus, the output lines of the shift register and the clock lines of the inverter are connected to each other through a parasitic capacitor.\nAs shown in FIG. 1, a ripple is generated in an anti-phase emission control signal EMB applied to output lines of a shift register because of an influence of emission shift clocks through a parasitic capacitor during a data holding period in a low-speed drive. The ripple reduces off-characteristics of some switches included in an inverter, and thus, drops a voltage level of an emission control signal EM output from the inverter. A reduction in the voltage level of the emission control signal EM changes an amount of electric current applied to OLEDs of the pixels, leading to a luminance drop."}
{"text": "1. Field of the Invention\nThe present general inventive concept relates to an apparatus to measure noise in an image signal and a method thereof, and more particularly, to an apparatus to measure an amount of noise in an image signal, usable with a process of attenuating the noise in the image signal, and a method thereof.\n2. Description of the Related Art\nNoise added to an image signal is the main cause of deterioration of image quality and lowering of performance of encoding and decoding of the image signal. Accordingly, many techniques for improving the image quality and the performance of encoding and decoding of the image signal by attenuating noise added to the image signal have been developed. In order to attenuate the noise added to the image signal, an amount of noise added to the image signal should first be measured accurately. For this, apparatuses for measuring noise in an image signal have been used. The amount of noise measured by the noise measurement apparatuses is generally expressed by a variance value or standard deviation that is the square root of the variance value.\nFIG. 1 is a block diagram illustrating a conventional apparatus for measuring noise in an image signal. Referring to FIG. 1, the conventional apparatus for measuring the noise in the image signal includes a SAD (Sum of Absolute Difference) calculation unit 10, a SAD comparison unit 20, a SAD counter 30, a comparison unit 40 and a variation counter 50.\nThe SAD calculation unit 10 calculates the sum of absolute difference (SAD) values among adjacent pixels of an input image. The SAD comparison unit 20 compares the calculated SAD value and detects whether each SAD value calculated by the SAD calculation unit 10 belongs to a predetermined range as a result of the comparison. The SAD comparison unit 20 outputs a true value if the calculated SAD value belongs to the predetermined range and outputs a false value if the calculated SAD value does not belong to the predetermined range. A counted value of the SAD counter 30 is increased if the output of the SAD comparison unit 20 is the true value. The SAD counter 30 is reset for a period of a picture (i.e., field or frame). Accordingly, the SAD counter 30 counts the SAD values that belong to the predetermined range for each picture.\nThe comparison unit 40 compares the value of the SAD counter 30 with a predetermined value. If the output value of the SAD counter 30 is less than the predetermined value, the comparison unit 40 makes the variation counter 50 perform an up-counting, while if the output value of the SAD counter 30 is greater than the predetermined value, the comparison unit 40 makes the variation counter 50 perform a down-counting, for the period of a picture. The output value of the variation counter 50, which is increased and decreased according to the comparison unit 40, becomes a noise measurement value, and this noise measurement value is fed back to the SAD comparison unit 20 to be referred to in determining the predetermined range.\nHowever, the conventional apparatus for measuring noise in the image signal as described above has the drawbacks in that the noise measurement value is calculated according to the SAD distribution of the image signal, and thus the noise measurement value may be differently calculated according to the characteristics of the image signals irrespective of the amount of noise applied to the image signals.\nFor example, even if the amount of noise applied to an image that is complicated or has many minute parts is the same as that applied to an image that is simple or has many flat parts, the noise measurement value may differently be calculated according to the SAD value distribution. Accordingly, it is necessary to calculate an accurate noise measurement value irrespective of the characteristics of the image signal."}
{"text": "Ovarian cancer is a significant health problem for women in the United States and throughout the world. Although advances have been made in detection and therapy of this cancer, no vaccine or other universally successful method for prevention or treatment is currently available. Management of the disease currently relies on a combination of early diagnosis and aggressive treatment, which may include one or more of a variety of treatments such as surgery, radiotherapy, chemotherapy and hormone therapy. The course of treatment for a particular cancer is often selected based on a variety of prognostic parameters, including an analysis of specific tumor markers. However, the use of established markers often leads to a result that is difficult to interpret, and high mortality continues to be observed in many cancer patients.\nImmunotherapies have the potential to substantially improve cancer treatment and survival. Such therapies may involve the generation or enhancement of an immune response to an ovarian carcinoma antigen. However, to date, relatively few ovarian carcinoma antigens are known and the generation of an immune response against such antigens has not been shown to be therapeutically beneficial.\nAccordingly, there is a need in the art for improved methods for identifying ovarian tumor antigens and for using such antigens in the therapy of ovarian cancer. The present invention fulfills these needs and further provides other related advantages."}
{"text": "The invention relates to an intermediate plate for mounting between housing parts of a fluid-operated control unit, in particular a gear unit, including a central unit that is constructed in the shape of a plate and, on either side of the central unit, a respective sealing system for sealing between the central unit and the housing part opposite the latter, this sealing system including sealing elements which are arranged on the respective side of the central unit and provide sealing around passages.\nIntermediate plates of this kind are known from the prior art.\nIn these, the problem arises of achieving the best possible seal between the central unit and the respective housing part that also maintains leak-tightness in the event of any kind of deformation of the housing parts and the central unit, even in the event of the high fluid pressures that are conventionally present in the case of a gear unit."}
{"text": "This invention relates to a cutting tool guiding system for metal cutting machines.\nIn the manufacture of steel plates and sheets, it is desirable to start with metal slabs which are an optimum width for the particular end product desired. Many such slabs are formed on continuous casting machines and are often wider than necessary. It is usually the practice to slit such slabs longitudinally into two or more strips.\nThe original slabs are often separated from the continuous cast strand by a flame cut-off device which utilizes a cutting torch which is moved transversely across the strand. The separated slab has often been found to have corners which are deformed in that they contain slag, or are rounded. These deformities have caused problems with prior guidance systems for slitting the slab when such systems followed the slab edge. Furthermore, at least some prior edge following devices have also had to contend with warped or misplaced slabs.\nThe present invention overcomes the aforementioned problems and provides a longitudinal cutting device which is edge-guided by a system which is not affected by corner deformities or warpage or misplacement of the slab.\nIn accordance with the various aspects of the invention, a cantilevered transverse arm is mounted for movement along both a longitudinal and transverse axis. The outer end of the arm is adapted to be disposed over a metal slab to be cut and carries a tool assembly which is adjustable thereon. A guiding assembly is fixedly mounted to the arm and includes a pair of longitudinally spaced guide rollers positioned so that the cutting tool is disposed in a transverse plane extending between the rollers. Each roller is individually shiftable into and out of engagement with the slab edge.\nInitially, the leading roller is advanced toward the slab edge, while the trailing roller is retracted. The rollers are then moved toward the slab edge until the leading roller engages the said edge ahead of the slab corner. As the entire device moves longitudinally along the slab, the leading roller is retracted while the trailing roller moves into engagement with the slab edge. Deformed slab corners are thereby bypassed by the tool guidance system."}
{"text": "In general, in a method for growing a silicon single crystal according to the Czochralski method, a raw material (polycrystalline silicon) is loaded inside a quartz crucible, polycrystalline silicon is melted by heat radiated from a heater to form a silicon melt, and then a silicon single crystal is grown from a surface of the silicon melt.\nA silicon wafer used as a substrate for a semiconductor device is added with a P-type dopant or an N-type dopant in the silicon single crystal growth process so as to have an appropriate resistivity value. In addition, the P-type dopant or the N-type dopant is further classified into a high-melting-point dopant whose melting point is higher than that of silicon and a low-melting-point dopant whose melting point is lower than that of silicon, and a method of adding a dopant to a silicon melt varies depending on a type of a dopant.\nA maximum resistivity value of an ingot, which may be realized with impurities included in a raw material (polysilicon, quartz crucible) currently used, is about 5 kΩ on the basis of the P-type. It is necessary to control a level of donor and acceptor included in a raw material in order to secure a resistivity of 8 kΩ or more.\nIn the related art, since impurities included in a raw material may not be accurately checked, a single crystal ingot was grown without introducing a dopant, and a part of a grown initial body was processed into a sample to measure a resistivity. In the case of this, since a procedure for calculating and inputting a necessary amount of dopant so as to have a target resistivity value is required and a growing time of a single crystal body for manufacturing a sample, an evaluation time of a sample, and doping time are added, there are problems that a process loss time is generated and the productivity of single crystal growth is also deteriorated."}
{"text": "1. Field of the Invention\nThe present invention relates to a cheese product containing a non-fat dry milk substitute product, and more particularly relates to a non-fat dry milk substitute product comprising lactalbumin, a modified whey solids product and sodium hexametaphosphate.\nAt the present time, there is considerable interest in the more efficient utilization of animal protein, including the recovery of animal protein from products formally considered to be waste materials of little potential use. Specifically, in the manufacture of cheese from milk, approximately half of the milk solids are coagulated as cheese, the remaining solids being contained in the residue or whey. Several processes have been developed for the recovery of valuable protein from whey and the separation of undesirable whey constituents from protein. Processes for treating whey are disclosed in U.S. Pat. No. 2,606,181, U.S. Pat. No. 3,002,823, U.S. Pat. No. 3,060,219, U.S. Pat. No. 3,547,900 and its Reissue U.S. Pat. No. Re. 27,806, U.S. Pat. No. 3,560,219 and U.S. Pat. No. 4,036,999. The entire disclosures of all of these patents are incorporated herein by reference.\n2. Description of the Prior Art\nAccompanying these developments in the processing of whey is the development and formulation of numerous products which contain whey and whey products. Specifically, whey and whey products are being used as non-fat dry milk substitute products for the non-fat dry milk normally used in cheese products.\nFor example, the use of a modified whey solids product in cheese products, in particular pasteurized process cheese food and pasteurized process cheese spread, as a non-fat dry milk substitute product is fully described in the assignee's copending U.S. patent application Ser. No. 680,613 filed on Apr. 26, 1976 and Ser. No. 818,645 filed on July 25, 1977, and concurrently filed Ser. No. 840,683, filed on Nov. 11, 1977 (Stauffer Case C-5072), the entire disclosures of which are herein incorporated by reference.\nThere is, however, still a need for a low cost non-fat milk substitute product for use in cheese products. Such a product would considerably lighten the economic burden to both the consumer and the manufacturer. While the United States consumption of American pasteurized process cheese (which does not contain non-fat dry milk) increased only 9% between 1969 and 1970, the consumption of pasteurized process cheese foods and pasteurized process cheese spread (which do contain non-fat dry milk) increased over 60% in 1970. Additionally, in the last few years, the price of non-fat dry milk more than tripled to about 65.cent. per pound in 1976.\nWhey and whey products, due to their low-cost and availability are a logical source for such non-fat dry milk substitute products. Difficulties arise, however, in formulating, from the numerous whey products which can be produced, a non-fat dry milk substitute product which is not only low in cost, but also does not adversely effect the properties of the cheese product.\nIt is thus an object of this invention to provide a non-fat dry milk substitute product for use in cheese products which is low in cost and does not adversely effect the properties of the cheese product."}
{"text": "JP 2014-163237 A discloses a conventionally known exhaust apparatus mounted on a vehicle with an engine. The exhaust apparatus disclosed in JP 2014-163237 A includes an aftertreatment device configured to purify exhaust gas emitted from an engine, an exhaust pipe through which the exhaust gas treated by the aftertreatment device is released into the atmosphere, and a nitrogen oxide (NOx) sensor mounted on the exhaust pipe. The NOx sensor detects a NOx concentration in the exhaust gas to be released into the atmosphere. The detection results are utilized for operation control of the aftertreatment device (for example, control of injection quantity of urea water for reduction of NOx) or utilized as information for diagnosing failure or deterioration of the aftertreatment device.\nHowever, NOx sensors are typically weak to water. For example, in a widely used NOx sensor equipped with a sensing portion containing ceramics such as zirconia, adhesion of water to a surface of the ceramics may result in a crack in the worst case, because evaporation of the water on the surface of the ceramics generates latent heat, thereby generating a possibility of large thermal stress in the ceramics. Moreover, since the NOx sensor is typically provided on an exhaust pipe that is communicated with the atmosphere, rain water or high-pressure cleaning water, for example, may enter the exhaust pipe and reach the NOx sensor. In view of the above, it is an important issue to protect the NOx sensor from moisture.\nIn this regard, the exhaust pipe disclosed in JP 2014-163237 A can effectively protect the NOx sensor from adhesion of external water such as rain water, because the exhaust pipe is formed in a U-shape with an intermediate part between upstream and downstream ends thereof, the intermediate part being the lowest part, and the NOx sensor is provided on the downstream end which is positioned higher than the intermediate part. The exhaust pipe, however, has to have a long passage length and a complicated shape in order to locate the intermediate part at a lower position than the other parts, thus involving significant increase in the size and in the cost of the apparatus."}
{"text": "Electromagnetic clutches are well known in the art and have been utilized commercially in many applications, including automobiles, for a number of years.\nA typical electromagnetic clutch includes a rotor which generally comprises an inner annular bearing portion, a clutch portion extending generally radially outwardly from one end of the inner portion, and an outer annular portion extending from the clutch portion in a generally overlying spaced relation with respect to the inner portion. The spacing between the inner and outer annular portions receives an electromagnetic coil to be energized to create a flux field in the rotor which allows attracting a floating plate mounted to the outer annular portion for selective coupling therewith when the electromagnetic coil is energized.\nThe clutch portion often includes a series of arcuate slots. The purpose of the arcuate slots is to direct the flux field back and forth between the floating plate and the clutch portion of the rotor for efficient magnetic attraction between the floating plate and the rotor.\nThe reverse clutch, as described in this patent, comprises an additional spring mechanism (metal or elastomeric, with or without pins) that allows for a floating plate to transfer torque to the hub plate by means of friction when the coil is de-activated.\nSuch a clutch allows selectively coupling and uncoupling of a rotatable shaft mounted to the hub plate."}
{"text": "One way that integrated circuit designers make faster and smaller integrated circuits is by reducing the separation distance between the individual elements that comprise the integrated circuit. This process of increasing the density of circuit elements across a substrate is typically referred to as increasing the level of device integration. In the process of designing integrated circuits with higher levels of integration, improved device constructions and fabrication methods have been developed.\nAn example of a common integrated circuit element is a transistor. Transistors are used in many different types of integrated circuits, including memory devices and processors. A typical transistor comprises a source, a drain, and a gate formed at the substrate surface. Recently, vertical transistor constructions that consume less substrate “real estate”, and thus that facilitate increasing the level of device integration, have been developed. Examples of vertical transistor constructions are disclosed in U.S. patent application Ser. No. 10/933,062, now U.S. Pat. No. 7,442,976, the entire disclosure of which is hereby incorporated by reference herein. While these improved transistor constructions are smaller and are packed more densely, they also often involve fabrication processes that are significantly more complex, therefore increasing fabrication time and expense. Fabrication complexity is increased even further when high density vertical transistors are formed in an array on the same substrate as logic circuitry that is positioned adjacent to the transistor array. In particular, conventional fabrication techniques use separate masks to independently define features in the device array region and in the device periphery region, since different process steps and materials are used to define the devices of these two regions.\nConventional semiconductor-based electronic storage devices, such as dynamic random access memory (“DRAM”) devices, include large numbers of transistor and capacitor elements that are grouped into memory cells. The memory cells that comprise a DRAM device are arranged into larger memory arrays that often comprise thousands, if not millions, of individual memory cells. Therefore, there is a continuing effort to reduce the complexity of the processes used to form densely-packed integrated circuit elements such as vertical transistor constructions."}
{"text": "Masonry construction techniques are well-known in the art. Such techniques, however, are impractical for many building projects due to the high cost of labor and materials. Additionally, because mortar must be applied by hand to each brick and the brick must be properly aligned with the bricks already in place, a high degree of skill is required to provide an aesthetically-pleasing result. Such results are extremely difficult to achieve with unskilled labor.\nThere have been attempts in the prior art to provide systems for assisting a laborer to properly aligning the bricks of a wall construction. Such systems are described in U.S. Pat. Nos. 2,172,816 and 3,170,267 to Douglas et al and Rosenfeld, respectively. In Douglas et al, for example, a plurality of T-shaped dowel plates are supported in aligned grooves along at least three edges of each brick. Rosenfeld describes a system wherein rod-like members are provided for vertical alignment with cross bars between such members for horizontal alignment.\nWhile such systems do provide improvements over manual techniques for aligning bricks in a wall construction, they do not adequately solve the problems associated with the use of unskilled labor. Moreover, the complexity of such systems make them difficult to use in practice and often more expensive than the cost of labor and materials. Such systems also do not reliably stabilize the bricks against lateral movement.\nThere is therefore a need to provide an improved wall construction technique that overcomes these and other problems associated with the prior art."}
{"text": "The present invention relates generally to the field of touch screen devices, and more particularly to selecting multiple areas of content on a touch screen device.\nA touch screen is an input device normally layered on the top of an electronic visual display of an information processing system. A user can give input or control the information processing system through simple or multi-touch gestures by touching the screen with a special stylus/pen and/or one or more fingers. Some touch screens use ordinary or specially coated gloves to work while others use a special stylus/pen only. The user can use the touch screen to react to what is displayed and to control how it is displayed (for example by zooming the text size). The touch screen enables the user to interact directly with what is displayed, rather than using a mouse, touchpad, or any other intermediate device (other than a stylus, which is optional for most modern touch screens).\nTouch screens are common in devices such as game consoles, personal computers, tablet computers, and smartphones. They may also be attached to computers or, as terminals, to networks. They also play a prominent role in the design of digital appliances such as personal digital assistants (PDAs), global positioning system (GPS) navigation devices, mobile phones, video games and some books (electronic books or E-books). The popularity of smartphones, tablets, and many types of information appliances is driving the demand and acceptance of common touch screens for portable and functional electronics. Touch screens are found in the medical field and in heavy industry, as well as for automated teller machines (ATMs), and kiosks such as museum displays or room automation, where keyboard and mouse systems do not allow a suitably intuitive, rapid, or accurate interaction by the user with the display's content."}
{"text": "Technical Field\nAspects of the present disclosure relate to a drying device and a printing apparatus.\nRelated Art\nA drying device is known that circulates a plurality of holders to hold members to be dried (drying target members) around an imaginary axis extending in a horizontal direction to dry at least drying target portions of the drying target members held by the holders."}
{"text": "This invention relates to improvements in photorefractive keratectomy (PRK).\nIn PRK, corneal tissue is removed in a controlled fashion to shape the surface of the cornea of a patient's eye to treat, e.g., myopia, hyperopia, presbyopia or astigmatism\nThe cornea comprises transparent avascular tissue that forms the anterior portion of the eye. The cornea functions as both a protective membrane and a \"window\" through which light passes as it proceeds to the retina. The transparency of the cornea is due to its uniform structure, avascularity, and deturgescence, which is the state of relative hydration of the corneal tissue.\nA major proportion of the refractive power of the eye is determined by the curvature of the anterior surface of the cornea, so that changing the shape of the cornea offers a way to significantly reduce a refractive problem in the eye.\nVarious techniques have been proposed for shaping the cornea of a patient's eye. Some techniques include removing the epithelium, and then shaping the underlying Bowman's and stroma layers. In PRK, photoablation is employed using e.g., ultraviolet radiation from an excimer laser, e.g., at 193 nm wavelength, or infrared laser radiation that has a wavelength in the range of about 2.9 to 3.2 .mu.m.\nIn one technique, described in Marshall et al., U.S. Pat. No. 4,941,093 (assigned to Summit Technology Inc.), the shape and size of the area of the corneal surface which is irradiated by laser radiation is selected and controlled so that some areas of the corneal surface become more eroded than others and a desired corneal shape is achieved.\nAnother technique, described in Muller, U.S. Pat. No. 4,856,513 (assigned to Summit Technology Inc.), uses a laser and an erodible mask. The mask, with a predefined profile of resistance to erosion by laser radiation, is disposed between the laser and the corneal surface. A portion of the laser radiation is absorbed by the mask, while another portion is transmitted to the corneal surface in accordance with the mask profile, thereby enabling the selective photoablation of the corneal surface into a desired shape.\nThere are circumstances in which it is desired to accomplish the PRK treatment with the ablated zone larger than 5 mm in diameter. With such zones, under usual operating conditions it has been observed that the final surface achieved by the ablation process differs from the expected shape. Surface irregularity and significant refractive error have been observed post-operatively in the corneal topographies of some patients treated for PRK. These irregularities may lead to visual disturbances such as diplopia, blurred vision, and loss of Best Corrected Visual Acuity (i.e., the vision obtained with the best possible lens correction).\nThe PRK procedure has achieved a clinically accepted level. However the possibility of achieving even better results with larger ablation zones has been somewhat elusive. The present invention provides a new insight into conditions that can occur in PRK and provides techniques that address these conditions to enable enhanced predictability, stability, and safety of the procedure to be achieved."}
{"text": "U.S. Pat. No. 3,122,137 describes a iontophoresis applicator applied on the orbit and not on the eye's surface and incorporating the current source. It does not propose means for keeping the eye's open and it is believed a large part of the product is also delivered into systemic circulation, due to a lack of precision in the placement of the device.\nIontophoresis devices in U.S. Pat. No. 5,522,864 and U.S. Pat. No. 6,101,411 should suffer from the same drawbacks.\nU.S. Pat. No. 4,564,016 discloses a device having a small application surface (diameter 1 mm), applied on sclera and allowing very high current densities (between 50 and 2000 mA/cm2) for “focal iontophoresis”. These values should be toxic for concerned tissues, as confirmed in Maurice article in Ophthalmology (January 1986, vol 93, number 1) entitled <<Iontophoresis of fluorescing into the posterior segment of the rabbit eye>>.\nLast, U.S. Pat. No. 6,154,671, discloses the principle of a device for delivering all kinds of active substances with safety and accuracy by iontophoresis, and answers then to most of the problematic of iontophoresis for ophthalmology.\nMore recently, U.S. Pat. No. 6,319,240 proposes an improvement of previous methods with a sealed reservoir applied on sclera (with a semi-permeable membrane on application face) under the eyelid. The semi-permeable membrane of this device is supposed to limit arcing effect between the electrode and the eye's surface that could occur due to the limited thickness and small surface of the device.\nU.S. Pat. No. 6,442,423 describes a device where the product is cast in a gel and applied on cornea.\nAll the previous devices are implemented for delivering substances through the cornea or the sclera, in order to reach some determinate tissues.\nNevertheless, there is a need to improve the efficiency and decrease the time of application of some therapies implemented by iontophoresis.\nFor example, the therapy of glaucoma by iontophoresis needs to be improved.\nGlaucoma is characterized by elevated intraocular pressure (IOP also known as ocular hypertension).\nGlaucoma is classified as “open-angle glaucoma”, which results from decreased permeability of the aqueous humor through the trabecular meshwork, or as “angle-closure glaucoma” that results from shifting of the iris forward so that the anterior chamber angle is obstructed or congenital.\nThe congenital form of glaucoma rarely responds to therapy and is more commonly treated with surgery, but not by iontophoresis means.\nIt is known that elevated IOP can be at least partially controlled by administering drugs which either reduce the production of aqueous humor within the eye, such as beta-blockers or carbonic anhydrase inhibitors, or increase the outflow of aqueous humor from the eye, such as miotics or sympathomimetics.\nThese pharmacological approaches help restore the IOP to a normotensive state either by inhibiting the production of aqueous humor by the ciliary body, or facilitating trabecular or uveoscleral aqueous humor outflow.\nTo this aim, most of the drugs are administered topically (eye drops) to avoid associated systemic effects.\nAlthough a wide variety of pharmaceutical treatments for lowering IOP are available for the glaucoma patient, these treatments are then limited either in terms of efficiency or side-effects.\nAdditionally, the previous iontophoresis devices are not specifically arranged for treating glaucoma.\nA first object of the invention is to provide an ocular iontophoresis device which leads to an increase of the concentration of drugs delivered into the intraocular tissue, improving for example the treatment of glaucoma.\nAnother object of the invention is to reach the first object by providing a iontophoresis device arranged for decreasing the loss of medication before reaching the targeted tissues by limiting the risks of absorptions of the medication by intermediates and/or by increasing the guiding of the ionized medication, for then more sufficiently treating intraocular tissues.\nAnother object of the invention is to decrease the time necessary for administrating a determinate amount of substances while administering the same amount of active substances."}
{"text": "The present invention relates to a method and apparatus for controlling the torque of an internal combustion engine.\nIt is known that there can occur conditions which result in the engine speed of a vehicle oscillating which in turn causes an oscillation in the torque output from the engine. The torque oscillation can sometimes be audible or even result in a physical oscillation of the vehicle which is commonly called \"kangarooing\" in the United Kingdom. In some vehicles, torque oscillation of the engine occurs due to a mismatch of the resonant frequencies of the engine, transmission and final drive and this is inherent at certain vehicle and engine speeds. It has already been proposed to mitigate the effects of this torque oscillation by monitoring the rate of change of engine speed and to alter the ignition timing so as to alter the torque output of the engine. In other words, when a drop in engine speed is detected by the monitoring operation, engine torque is increased by advancing the ignition in order to counteract the drop. An increase in engine speed is likewise counteracted by retarding the ignition so as to damp out the oscillation.\nThis system is only possible because the resonant conditions are known and also the amount of phase shift between engine and wheel speed of the vehicle is a constant under these conditions. The requisite ignition timing variation can therefore be evaluated.\nThere exists another set of conditions which may induce \"kangarooing\" and that is during acceleration. Here, the phase shift between engine and wheel speed is not constant and so the technique proposed above cannot be used."}
{"text": "In mobile phone applications, for example, an antenna may be coupled to a power amplifier (PA) to drive the antenna circuit. To drive the antenna circuit effectively, the output of the PA should be impedance matched to the antenna circuit. One method involves the use of inductors for impedance matching of the output. However, a PA output match inductor forms a strong magnetic field that couples to adjacent devices, such as surface mount device (SMD) parts and duplexers. This magnetic coupling contaminates receive (RX) and transmit (TX) signals from antennas driven by the PA, which worsens the sensitivity of the PA/antenna combination and noise figure for the amplifier/antenna circuit. If a conformal shield is used to shield the entire device—PA, antenna, SMDs, substrate, etc. (module level conformal shield)—the coupling becomes worse since the magnetic waves will bounce around inside the module and the shield itself will work as an antenna.\nAccordingly, there is a need for systems, apparatus, and methods that overcome the deficiencies of conventional approaches including the methods, system and apparatus provided hereby."}
{"text": "This application relates generally to payment accounts. More specifically, this application relates to systems and methods for activating a payment account.\nGift cards, stored value payment instruments, and other prepaid payment instruments have become increasingly popular in recent years. These gifts cards or payment instruments often are linked to an underlying payment account with an account balance, although the underlying payment account may not be tied to a particular individual. Each of these payment instruments also often has a number that is encoded in a magnetic stripe and/or a bar code, with the number being associated with the underlying payment account.\nThe increasing popularity of these payment instruments have given rise to several problems. One problem of great concern is security. In one scenario, a thief may go to a store and secretly scan the bar code of a payment instrument that has not been activated. The thief would then simply wait for an unknowing purchaser to buy the payment instrument. Once the payment instrument has been activated following the purchase, the thief may clone the number onto another “evil twin” payment instrument, and either resell or use the “evil twin” payment instrument as if the thief was the legitimate purchaser.\nUnlike traditional magnetic stripes and bar codes, radio frequency devices can provide greater security because they are better able to prevent a number that is encoded in them from being read and/or cloned. In some instances, the number that is encoded in a radio frequency device may change each time the radio frequency device is read. However, if radio frequency devices are used, it may not be possible to activate the payment account readily because many point of sale devices remain incapable of reading radio frequency devices.\nHence, there is a great need in the art for systems and methods to improve activation of a payment account."}
{"text": "The embodiment of present invention relates to a film stripping mechanism.\nIn the display field, substrates, such as glass, plastics and the like, are required when making a display device, and the surface of the substrate is attached with films when being rolled out so as to ensure the integrity and cleanliness of the substrate, but during subsequent processes of using the substrate, the film thereon should be stripped away, and the prior art commonly takes off the film on the substrate manually in a manner of scraping using a knife. However this method has a greater force against the substrate, tends to cause damage to the substrate, and the working efficiency thereof is low, such that the operator may also bring about negative effect upon the substrate if he/she lacks of skill, causing waste of materials."}
{"text": "The prior art is replete with innovations in contact lens designs. One such innovation is the rigid gas permeable (RGP) contact lens. This optical device has been successfully used in corneal lens to provide visual corrections of hyperopia, presbyopia, and myopia. The lenses have also been used for limited corneal molding. Examples of such applications are described in U.S. Pat. Nos. 4,952,045, 5,192,365 and 5,349,395 issued to Nick Stoyan. Many disadvantages still exist, however, with conventional contact lens designs. Corneal lens, for example, are generally prone to dislodgement from the eye if the wearer is subjected to sharp changes in motion, as often encountered during participation in sports and other physical activities. Dislodging of the lens can be both uncomfortable and dangerous.\nAnother type of lens is the scleral contact lens. A scleral lens covers most of the sclera and all of the cornea of the eye. Additionally, the scleral lens is wholly supported by the scleral portion of the eye and vaults the cornea to provide a cavity for tears or a tear substitute. These lenses have been conventionally formed of blown glass and have been relatively thick, being on the order of 2.0 to 3.0 mm thick. The concave, or posterior, surface of the lenses had two primary areas of curvature, the first being an outer peripheral area that approximated the scleral curvature, and the other being a central curved area that vaulted the cornea by as much as 2.0 mm providing a chamber that was filled with a tear substitute. The convex, or anterior, surface of the lens had a contour approximating that of the inner surface except in the center which was ground to provide optical correction. In addition to being thick, scleral lenses are also relatively large, having a typical diameter of about 25 mm, often causing eye irritation if worn for more than several hours.\nPrior art scleral contact lenses have also been manufactured of plastic materials based on a casting of the eye. In this arrangement, the center of the lens has minimal vaulting and does not require a tear substitute. However, the diameter and thickness of the lens is virtually the same as the above mentioned glass, and similar glass-plastic composite lenses.\nPlastic scleral contact lenses were also produced by lathing and polishing a plastic preform. The concave surface contained at least two, and as many as 10, curves to approximate the shape of the eye. These designs produced a relatively thinner lens (albeit still very thick by today's standards) with little or no clearance in the scleral area and little to extreme clearance in the corneal area. While some of the above designs were considered improvements over earlier lenses, the ability to wear scleral contact lenses has, heretofore, been very limited, both as to wearing time and in the number of patients able to tolerate them.\nOne attempt to provide a smaller diameter contact lens is described in U.S. Pat. No. 4,194,815 issued Mar. 25, 1980 to Wayne E. Trombley for a SEMI-SCLERAL CONTACT LENS. The Trombley lens is aligned with, i.e., contacts, the central corneal area of an eye, then vaults the limbus and a portion of the scleral area of the eye before coming down and making contact, at the very outer peripheral area of the lens, with the sclera. Thus, the Trombley lens vaults all of the eye area between the central cornea and a point contact line on the sclera. If formed of a rigid material, a possibility suggested by Trombley, the very small bearing surface on the sclera would be very uncomfortable and would permit the lens to move, particularly in the up and down directions, creating displacement of the lens' curvature at the center. In order to avoid adverse movement of the lens in view of the small scleral bearing area, the teaching in Trombley is obviously only applicable to a pliable lens in which the sagittal depth is not critical.\nThe present invention is directed to solving the problems set forth above. It is therefore desirable to have a contact lens that conforms with the limbus and a small portion of the sclera to provide a distributed outer bearing area for the lens, and extends under a portion of the eyelid so that the lens is not easily dislodged, such as is often the case with corneal contact lenses. It is also desirable to have a contact lens that is relatively thin, comfortable to wear, and easy to insert and remove. Furthermore it is desirable to have a contact lens that can be worn after corneal surgery, such as keratotomy, in cases where the surgical procedure was not completely successful in achieving the desired visual acuity. Still further, it is desirable to have a contact lens that can mold the cornea to a desired shape by controllably distributing the bearing pressure of the eyelid on both the corneal and scleral portions of the eye."}
{"text": "The present invention relates to protective eye wear, and more particularly, to goggles having transparencies which cover the lens and are torn away during a race instead of cleaning the lens.\nMotorcross racers typically drive their motocycles at high speeds across rough terrain. Along with various other equipment for protecting the body, these racers typically wear goggles. During a race, dirt and mud build up on the lens of goggles. This substantially impairs the racer's vision and if not corrected, can cause an accident and serious injury. However, because the racer is required to keep both of his hands on the respective handlebar grips of the motorcycle to maintain control, it is extremely difficult to effectively wipe the lens clean during a race. The amount of time that one of the racer's hands would have to be off the handlebar grip could result in a loss of control and an accident. If the racer slows down sufficiently to safely clean the lens, then it is likely that he will be passed by other racers and will loose the race.\nOne solution to this problem has been provided by a commercially available goggle in which a webbing of clean transparent material is fed across the lens of the goggle between rollers mounted at the opposite side edges of the goggles. When the segment of transparent webbing overlying the lens becomes too dirty, the racer pulls on a cord and mechanisms cause a clean segment of the transparent webbing to be pulled across the lens. However, these goggles require mechanical feed mechanisms which are subject to failure and the complexity of this solution makes it generally one that has not yet received wide spread receptance.\nAnother solution to this problem of dirty goggles that has received widespread acceptance is the use of tear off transparencies. These comprise individual, thin pieces of transparent plastic material having a general shape of the goggle lens. The transparencies normally have tabs around their peripheral edges. A plurality of these transparencies may be overlaid on the goggle lens, and their tabs inserted between the lens and the flexible lens frame to hold them into position. Generally, these transparencies have an extended arm portion at one side thereof which extends past the side edge of the goggle. These arm portions may extend at different angles in the general plane of the lens. During a race, dirt and mud builds up on the outermost transparency. When the vision is sufficiently obscured, the motorcross racer can reach up with one hand and tear away one of the transparencies with one quick motion. The dirty transparency then falls to the ground and the racer has a clear view through the lens and remaining unsoiled transparencies. A principal drawback in the foregoing arrangement is that the transparencies are held to the lens only by the tabs. Furthermore, the number of transparencies that can be used is limited to a small number, for example three, since there is only a limited amount of angle through which arm portions of the transparencies can extend, and not yet overlap.\nRecently, a new type of goggles equipped with tear off transparencies has been marketed under the trademark OAKLEY. In this goggle, an element with a small round head projects outwardly from the lens adjacent one side edge thereof. The transparencies are all identical, each having an arm portion which extends from the left side edge thereof viewed from the racer's standpoint. The arm portions have a large circular hole adjacent their outer end. Each of these transparencies further has insertion tabs around the peripheral side edges thereof and a small hole which tightly fits over the rounded portion of the element which projects from the lens. When a plurality of these transparencies are inserted over the lens of the goggles, they are put in place one at a time.\nThe transparencies used with the OAKLEY goggles are inserted as follows. The tabs of a first transparency are inserted between the lens and the flexible lens frame. Then the portion of the transparency having the small hole in it is forced over the rounded end of the projecting element. The arm portion of the first transparency is then curled back on itself, away from the goggles. The outer ring of the arm portion is hooked around the projecting element. The arm portion of the first transparency is held in its folded configuration with its ring portion engaged with the element until the next transparency is put in place and the projecting element is forced through its small hold. This serves to retain the arm portion of the first transparency immediately below in its folded state. Each succeeding transparency is similarly installed on the goggle. The arm portion of the last transparency is not folded but is left hanging out, extending beyond the folded arm portions of the transparencies below.\nDuring a race, a motorcycle racer wearing the OAKLEY goggles can easily grasp the arm portion of the outermost transparency and pull this transparency free of the projecting element and from the lens with one quick motion. When this is done, the arm portion of the transparency immediately beneath unfolds, and is ready for grasping by the racer when it time once again to clean his goggles.\nWhile the arrangement utilized in the OAKLEY goggles has a number of advantages, its principal disadvantage is that a special lens must be fashioned with a projecting element which extends therethrough adjacent one side edge of the lens. Furthermore, the projecting element in the lens tends to obstruct the view since the projecting element is opaque. The folded arm portions of the multiple transparencies tend to be clumped over the side edge portion of the lens. It would be desirable to provide a means of retrofitting existing goggles so that they could use the folding type of transparencies without having to replace their lens."}
{"text": "The present invention relates to an image sensing apparatus and a recording/reproduction apparatus having a protection cover such as an electronic still camera and a video camera.\nRecently, there have been proposed a type of digital still cameras and video cameras having a lens-protection cover which is open to expose an objective lens during image-capturing and is closed to protect the lens during non image capturing. Such lens-protection cover is referred to as xe2x80x9cbarrierxe2x80x9d hereinafter.\nFIG. 1 is a perspective view showing a digital camera equipped with such barrier.\nIn FIG. 1, reference numeral 60 denotes an operation lever for opening and closing a barrier (not shown). A user operates this operation lever 60 to open the barrier connected thereto.\nWith the spread of personal computers, it has been desired to control a digital camera from a computer.\nIn the prior arts set forth above, an operation to open or close the barrier is made at the camera body side. The prior arts have not paid any attentions to how and when to open and close the barrier when the camera is controlled by the external computers,\nThe present invention has been made in consideration of the aforementioned problem involved in the prior art.\nAn object of the present invention is to provide an image sensing apparatus having a protection cover, such as a digital still camera and video camera, that can be controlled depending on whether an external controller such as a PC and workstation is connected to the apparatus.\nThe above object is achieved by providing an image sensing apparatus comprising: a protection cover which is openable and closable, and protects an image sensing optical system; a cover driving device driving the protection cover; and a determination device determining whether the image sensing apparatus is set into a image sensing mode, or into an external control mode in which the apparatus is controlled by an external controller unit, wherein the determination device determines an operation of the cover driving device in accordance with a determination result.\nAccording to the apparatus thus constructed, in a case where the image sensing apparatus is controlled by the external controller unit, opening and closing operations of the cover are suitably controlled.\nOther features and advantages of the present invention will be apparent from the following description taken in conjunction with the accompanying drawings, in which like reference characters designate the same or similar parts throughout the figures thereof."}
{"text": "1. Technical Field\nThe present disclosure relates to multiple-cylinder internal combustion engines having intake and/or exhaust valves operated by a camshaft positioned in an engine block with an associated valvetrain.\n2. Background Art\nConventional internal combustion engines use a camshaft-driven valvetrain to operate intake and exhaust valves that control the exchange of gases in the combustion chambers formed between the engine block and cylinder head. Engines are often categorized by the location of the camshaft relative to the valves, with overhead cam valvetrains driven by a camshaft in the cylinder head over the valves, and pushrod valvetrains or “cam-in-block” valvetrains having the camshaft located in the engine block with the valves operated using pushrods and rocker arms.\nCurrent four-valve-per-cylinder pushrod engines include two intake valves and two exhaust valves for each cylinder. Each pair of valves is operated in tandem by a bridged valvetrain that includes a camshaft-driven cam follower (also referred to as a tappet or lifter) connected by a single pushrod to a rocker arm that drives a bridge coupled to the pair of valves (intake or exhaust). This bridged valvetrain is a cost-efficient design that achieves acceptable performance for many applications, although operation of the two bridged valves is not precisely synchronized because the force exerted on the bridge can not be perfectly balanced between the valves, the valves may have slightly different spring forces, and the valve components may experience slightly different wear. This may result in one valve opening late and/or one valve may seat first while closing causing the other valve to seat late with a higher than intended velocity. In addition, valve stem tips are edge loaded by the bridge with higher stresses resulting in higher rates of wear and potential noise, vibration, and harshness (NVH) concerns. While single overhead cam (SOHC) and dual overhead cam (DOHC) systems have independently controlled valves to address some of these issues, the SOHC and DOHC systems are significantly more expensive and have large package width relative to a cam-in-block design.\nTo provide various advantages over conventional pushrod, SOHC, and DOHC engines, an engine and valvetrain having dual pushrod lifters and independent lash adjustment has been developed as described in commonly owned and copending U.S. patent application Ser. No. 11/164,620 filed Nov. 30, 2005. While suitable for many applications, the number of pushrods utilized may impose packaging constraints on port placement in the cylinder head."}
{"text": "Conventional methods of placing telephone calls utilizing so called “reverse billing”, wherein a called party assumes the cost of the call, involve the placing of collect telephone calls. A collect telephone call can be placed by a calling party and the called party may then accept or deny the call and the associated charges. The cost of the collect call is billed to the called number and appears on the telephone billing statement of the called number. Thus, the collect call billing method is restricted to the telephone service account of the called number, i.e., the specific destination number.\nIt is well known that penal facilities are largely filled with poor people. The families of these poor incarcerated persons are often poor also, and may have their phones prohibited (blocked) from receiving collect calls and/or disabled from making toll calls (e.g., “900” type calls). Incarcerated persons are almost exclusively limited to making collect calls. Thus, it can be expensive, if not impossible for an incarcerated person to call a family member. Further, incarcerated persons have limited opportunities to make phone calls due to restricted access to telephones. In addition, friends and family members cannot directly call an incarcerated person, as the incarcerated person does not have a phone in his or her individual cell. Therefore, there remains a need for an improved channel of communication between incarcerated persons and their families, so as to ease the financial and scheduling burdens normally placed on the families of incarcerated persons.\nOn the other hand, a calling party has the option of placing a telephone call from any telephone and paying for the call by alternate means, such as with a credit card or a prepaid calling card. In these situations, the calling party is usually the owner of the card or account. Thus, they are in control of the account and any associated Customer Identification and Verification (CIV), such as a Personal Identification Number (PIN).\nIn many situations, a called party, such as a relative, friend, or business associate of a calling party, may have an important need for establishing telephone communications with the calling party that allows the called party to accept the cost of the call while maintaining ownership and control of the charged account. For example, a parent may have an important need to allow a child to call home from a public telephone without using a calling card, credit card, collect charging, or coins. Another example exists in connection with correctional facilities, where family members may wish to receive calls from an inmate and accept the charges of the call without incurring higher costs associated with normal collect calling. Yet another example exist where the called party may be at a destination number where accepting charges to that account would be inappropriate, such as a parent or husband at their place of employment.\nU.S. Patent Application Pub. No. 2003/0198325 discloses specialized pre-paid calling and voice messaging systems and methods for allowing inmates to communicate with parties outside of custodial facilities. Calls to and from an inmate are pre-paid by way of advertising revenue and/or by way of pre-paid calling cards, smart cards, credit cards, etc. supplied by the outside parties. A call processing system may be configured to receive a call from an inmate and to connect the call to an outside party able to provide up-front payment for the call. If the outside party is not able to provide up-front payment for the call, the call processing system may allow the inmate to record a voice message that can subsequently be accessed by the outside party. An access fee may be charged to the outside party for access to the voice messaging system. The call processing center may also be configured to present targeted advertising messages to inmates and/or to outside parties at selected times.\nAnother system disclosed in U.S. Pat. No. 6,282,274 enables a telecommunications service subscriber to designate billing options for outgoing and incoming calls on a per-call basis, even if the subscriber's account includes several numbers and/or addresses and also preferably allows users to designate default billing accounts on a per-number basis. The subscriber may have a plurality of service accounts established in a database such as service profile. Calls for personal use and calls for other uses (e.g., business calls, separate client accounts, etc.) may be distinguished on a per-call basis by the call originator for outgoing calls or by the call recipient for incoming calls if the recipient is billed for incoming calls. For both outgoing and incoming calls, network usage is automatically allocated to the appropriate service account.\nU.S. Pat. No. 6,639,977 discloses a system and method that allows a calling party to call a destination number and, upon approval of the called party, reverse the billing of the telephone call so that it is deducted from an account owned and maintained by the called party other than the account associated with the destination number. The method includes the steps of receiving a predetermined access number from a calling party; prompting for a destination number; notifying a called party of the call they are receiving; allowing the called party to accept the call and enter a number associated with an account other than the telephone service account of the destination telephone number; verifying the account number and whether the account has sufficient value available; completing the telephone call to the destination number; and processing the charges associated with the call to be charged against the account.\nU.S. Pat. No. 6,668,045 discloses a computerized system that allows inmates of a correctional institution to communicate with parties outside the facilities via telephonic or electronic messaging events. The system provides a means to bill for the events while controlling to whom the inmate can communicate with or from whom they may receive communications, including controlling the method, the length, limited content editing, the frequency of communicating, storage of and the manner of alerting a recipient of a message or receipt thereof.\nHowever, no current methods are known that allow the called party the same flexibility that is to direct charges of a telephone call to a credit card, debit card, or prepaid calling card."}
{"text": "Apparatuses having movable mechanical parts are present in most industrial contexts. A typical such apparatus has some kind of motor or actuator acting on a body that is movable in a translational and/or rotational sense. The body may interact with other mechanical parts or components in order to perform one or more predetermined tasks. These tasks may be of various types, e.g. closing a valve, moving a part between two predetermined positions, applying a force, adjusting positions, compression of gases, pumping of fluids, rotating filters, fans etc.\nMonitoring of mechanical apparatuses for the purpose of detecting signs of maintenance need or erroneous operation is requested in many areas. Emergency faults typically involve large costs in terms of production losses, spare parts and short-notice labour costs. An early detection of e.g. wear or faults may allow any maintenance operations to be performed in a well-planned manner, and may reduce the actual damage caused in the apparatuses. A common approach is to supply the apparatuses with different kinds of sensors, measuring vibrations, torques, forces, temperatures, pressures, sound etc. The sensors are then connected to a monitoring or diagnosing equipment performing an evaluation of the measured quantities.\nIn many industrial applications, mechanical apparatuses are sometimes present in areas with hazardous or aggressive environments, in limited spaces or simply in positions which are difficult to reach from outside. In such cases, it may be difficult to provide sensors directly on the apparatus or connections to the evaluation devices. A remote monitoring or diagnosing is thus required.\nWhen mechanical equipment is re-mounted when being opened for inspection, there is a finite probability that the mechanical equipment is not in operative condition or its settings for proper operation have been changed. In purpose to avoid future need of maintenance to correct the introduced errors it is valuable to find such errors in advance.\nIn the U.S. Pat. No. 4,965,513 is a motor current signature analysis method disclosed. A current supplied to an electrical motor driving e.g. a valve is measured at a position remote from the actual motor. The current signal is analysed in the frequency domain, i.e. in principle a vibrational analysis. The frequency spectrum comprises typically a number of more or less pronounced peaks, which are possible to relate to mechanical parts performing a periodic motion. By comparing the amplitudes of the frequency components in a spectrum with a spectrum recorded when a fault-free operation was assured, trends may indicate abnormalities or degradation. There are some remaining problems with the method disclosed in U.S. Pat. No. 4,965,513. Non-periodic motions are not possible at all to monitor or evaluate and the evaluation is basically performed as an associative evaluation, i.e. comparison with earlier spectra, which means that even if some trends are detected, there is basically no detailed information about what fault signal really is present.\nAlso the monitoring system disclosed in WO 99/05501 is based on an analysis of cyclical tasks. The power consumption as a function of time is compared with pre-determined upper and lower limit thresholds. The number of crossings of the thresholds is counted and evaluations are performed based on these figures and according to pre-determined sample recordings for a set of known faults. Here, the system is trained for a specific task by measuring the operation at well-defined situations. The analysis is a pure associative procedure. This implies that the total system has to be rather simple, in order to be able to cover possible behaviours. Also, all faults that are requested to be identified have to be caused on purpose and measured in order to provide the associative sample recordings.\nThe U.S. Pat. No. 4,888,996 describes a method and a device for registration of the work performed when a DC motor operates a valve. The mechanical output torque is determined by registering the electrical power supplied to the motor and compensating for electrical losses in the motor. The output torque, or a quantity proportional thereto is monitored as a function of time. A series of events can be monitored and interpreted as different operational phases. Certain quantities related to the state of the valve and its operation are possible to diagnose by relating the main features of the output torque curve to the course of events occurring when operating the valve. However, the analysis is based on simple and well-known models related to dimensioning of actuators.\nThe analysis in U.S. Pat. No. 4,888,996 is based entirely on static interpretations of the valve as one collective item. Such analysis is typically based on differences between different operation phases, where different sets of movable parts are in operation. Referring to FIG. 6 in U.S. Pat. No. 4,888,996, “c” refers to a duration of a first operation phase, “a” refers to a difference in collective average load between a fourth and second operation phase, “b” again refers to a duration of an operation phase, in this case the fourth one and “d” refers to the total load at the onset of the sixth operation phase. All such parameters are closely linked to changes in operation phase and are based on a collective view of the system. Although some information about the actual existence of operation problems is possible to extract, it is not possible to conclude what the details causing the problem are.\nThe U.S. Pat. No. 5,748,469 discloses a method of detecting errors in components in an automatically operated valve. A mathematical model of a control valve assembly is defined. The mathematical model is fitted to data that consists at least several values of a control signal and corresponding values of a position of the control valve. The analysis is thus based on a response analysis between an input signal and a resulting behaviour of a measured quantity, in this case a position. For deeper analysis, further measured data are needed. Certain parameters of the mathematical model are selected as “critical” parameters. The “critical parameters” are supposed to be mainly dependent on a part system, which is why the occurrence of errors readily may be detected to a certain part system if the associated critical parameter changes considerably.\nThe method is entirely dependent on a measured position as well as a control signal as an input, which as discussed above is unacceptable in many applications. The model is further based on a collective approach, i.e. where the valve-operating device is considered as one unit having statistically averaged properties that can be expressed as parameters. This is emphasized in that e.g. one single friction coefficient is used that is only dependent on the rotational angle. Furthermore, a reduced mass is used for the piston, i.e. a weighted averaged value of different interacting masses. Also a combined inertia moment for the valve and piston is used. Such a model thus enables a description of the behaviour of the valve system as one integrated single item. However, even if such errors may be detected and associated to a certain part system, such models do not provide any means to efficiently precisely localise the error within the part, and certainly not to estimate any error sizes."}
{"text": "When a user of a messaging system, e.g., electronic mail, is out of the office for a period of time, the user typically activates an automated reply function. The automated reply function automatically sends a reply in response to each message that is received in the user's messaging account. The automated reply typically includes text indicating when the user will be back in the office and able to respond to the message that triggered the automated reply.\nIn some cases, the user is out of the office for longer than anticipated without access to the messaging system. As the date noted in the automated reply comes and goes, people attempting to contact the absent user that previously received an automated reply expect some sort of follow-up action to occur based upon the assumption that the absent user is back in the office.\nWhen no follow-up action occurs, people may become frustrated by the lack of action while remaining unaware that the absent user is still not in the office. Further, when one sends a further message to the absent user, the same automated reply may be received in response, which now includes incorrect and out-of-date information. Alternatively, no further automated message is sent. In either case, the person trying to contact the absent user may interpret the lack of action or contact that the absent user is choosing not to respond when in fact the absent user is unable to respond."}
{"text": "1. Field of the Invention\nThe present general inventive concept relates to a method of manufacturing an inkjet printhead, and an inkjet printhead manufactured using the method, and more particularly, to a method of manufacturing an inkjet printhead in which formation of a reaction layer at the bottom of a nozzle layer is prevented, and an inkjet printhead prepared using the method.\n2. Description of the Related Art\nInkjet printheads eject tiny droplets of print ink onto a predetermined portion of a to-be-printed sheet so as to produce a predetermined color image, and are categorized into thermal driving type inkjet printheads and piezoelectric driving type inkjet printheads according to an ejection mechanism of ink droplets. As for the thermal driving type inkjet printheads, ink droplets are ejected by an expansion force of bubbles generated when a heat source is applied to the generated bubbles within the ink. As for the piezoelectric driving type inkjet printheads, ink droplets are ejected by a pressure to ink due to a deformation of a piezoelectric device using the device. FIG. 1 is a cross-sectional view of a conventional thermal inkjet printhead.\nReferring to FIG. 1, the conventional thermal driving type inkjet printhead includes a substrate 10, a channel forming layer 20 formed on the substrate 10, and a nozzle layer 30 formed on the channel forming layer 20. The substrate 10 has an ink feed hole 51, and the channel forming layer 20 has an ink chamber 53 that can be filled with ink and a restrictor 52 connecting the ink chamber 53 to the ink feed hole 51. The nozzle layer 30 has at least one nozzle 54 through which ink is ejected from the ink chamber 53. On the substrate 10, at least one heater 41 is mounted to heat ink within the ink chamber 53 and at least one electrode 42 is mounted to supply a current to the heater 41.\nAn ink droplet ejection mechanism of the conventional thermal driving type inkjet printhead will now be described in detail. Ink is fed from an ink container (not illustrated) into the ink chamber 53 through the ink feed hole 51 and the restrictor 52. The ink filled within the ink chamber 53 is then heated by the heater 41 which is formed of a heat-generating resistance and which is located within the ink chamber 53. Once the ink boils, ink generates bubbles, and the generated ink bubbles apply pressure to the ink filled within the ink chamber 53. Therefore, the ink within the ink chamber 53 is ejected through the nozzle 54 to the outside of the ink chamber 53 in the form of droplets.\nA conventional inkjet printhead has a protrusion 55, hereinafter referred to as a reaction layer, at the bottom of a nozzle layer, as illustrated in FIGS. 3A and 3B. The reaction layer is formed when a photoresist that is used to form an ink channel is mixed or reacts with a photoresist composition that is used to form a nozzle layer, and degrades the ejection performance of the ink."}
{"text": "In a conventional high dynamic range (HDR) image synthesis method, multiple images that correspond to different exposure values (e.g., −1 EV, 0 EV and +1 EV) are acquired by autobracketing, and are then synthesized to generate an HDR image. However, in some photographing scenes, image quality of the HDR image obtained by the abovementioned method may not be good enough."}
{"text": "(Not Applicable)\n(Not Applicable)\nThe present invention relates generally to chip modules, and more particularly to a chip module including a highly configurable chip array which is interconnectable to a universal package capable of hermetically sealing the chip array.\nThere is currently known in the prior art chip modules comprising a chip array which is mounted into a package for purposes of protecting the chip array from contaminants such as dust and moisture. The package may also be used to convert the chip array to a particular standard pin format. The chip array typically comprises a substrate (e.g., a circuit board) having off-the-shelf, pre-packaged memory or logic devices (e.g., TSOP packaged chips) mounted to one or both of the opposed sides or faces thereof. This chip array is mounted into and sealed within the associated package of the chip module. These prior art chip modules wherein the chip array is sealed within the package to protect the same from contaminants are often used to achieve compliance with stringent military standards such as the JET xe2x80x9cXxe2x80x9d standard which require extremely low susceptibility to component failure as a result of contamination from external or environmental sources.\nThough prior art chip modules are generally suitable for protecting the chip array within the associated package from external contamination, they possess numerous deficiencies which detract from their overall utility. One such deficiency lies in the inability to test the currently known chip modules until the completion of the package assembly process. Due to the complexity of their construction which is attributable in large part to the heat and moisture sensitivity of the internal chip array, such chip modules experience a relatively high failure rate during initial testing. Another deficiency lies in the inability to recover components from these chip modules in the event the same are unsuccessfully tested. For example, if the failure of the chip module during testing is attributable to a soldering problem, the chip module must be scrapped in its entirety due to the currently known configurations and assembly methods related thereto not allowing for the recovery or salvage of either the chip array or the package.\nA further deficiency of the prior art chip modules is that the package thereof is device specific, i.e., tailored to accommodate a particular chip array. As such, since the chip array is not confined to one specific configuration but may have any one of a multitude of differing configurations, any particular chip array will typically require a different device specific package. In this respect, the development of a new chip array carries with it the need to develop a new package, and hence new tooling as well as a new assembly protocol. As will be recognized, these constantly changing tooling and assembly requirements significantly increase the product costs of prior art chip modules. Moreover, the configuration of such prior art chip modules, and in particular the chip arrays thereof, does not lend itself to use with newly developed memory or logic devices, thus increasing susceptibility to obsolescence.\nThe present invention addresses and overcomes these deficiencies of prior art chip modules by providing a low cost, high reliability, high density chip module wherein a highly configurable, custom chip array is interconnectable to a xe2x80x9cuniversalxe2x80x9d package. In the present chip module, the chip array includes a novel and unique interconnect substrate having an interconnect array which is specifically adapted to electrically communicate with a corresponding interconnect array provided within a cavity of the package. In view of the complementary nature of these interconnect arrays, the package is universal in the sense that it may be used in conjunction with any chip array including the interconnect substrate. The package in the present chip module, in addition to being usable with any chip array including the interconnect substrate, is adapted to convert such chip array to a standard sixty-eight (68) pin format.\nIn addition to providing economies in the manufacturing/assembly process and thereby significantly reducing product cost, the configuration of the present chip module and assembly methodology preferably employed in relation thereto allows for the testing of the chip array prior to the mounting of the same within the package, and the testing of the combined chip array/package prior to the chip array being sealed within the package. As a result, in the present chip module, either the chip array or the package may usually be salvaged in the event of an unsuccessful result being obtained during the testing process. Moreover, the configuration of the chip array of the present chip module provides a high level of adaptability to new memory or logic devices, thus substantially eliminating its susceptibility to rapid obsolescence. The completed chip module formed in accordance with the present invention also satisfies the most stringent military standards. These, and other advantages attributable to the present invention, will be described in more detail below.\nIn accordance with the present invention, there is provided a chip module which includes a universal package for accommodating any one of uniquely configured chip arrays (e.g., memory or logic arrays). In the preferred embodiment, the chip array, which is a sub-assembly of the chip module, comprises an interconnect substrate having opposed, generally planar surfaces and including a first interconnect pad array disposed on at least one of the surfaces thereof. The interconnect substrate preferably has a generally square configuration defining four peripheral edge segments, with the first interconnect pad array extending along one of the peripheral edge segments. The first interconnect pad array itself preferably comprises first and second sets of conductive interconnect pads which are disposed on respective ones of the opposed surfaces of the interconnect substrate and extend along a common peripheral edge segment thereof. The conductive interconnect pads of the first and second sets are arranged in identical patterns such that the conductive interconnect pads of the first set are aligned with respective ones of the conductive interconnect pads of the second set. Additionally, a plurality of vias are preferably extended through the interconnect substrate between respective pairs of the conductive interconnect pads of the first and second sets thus establishing electrical contact or communication between the pads of each pair.\nIn addition to the interconnect substrate, the chip array of the present chip module comprises at least one integrated circuit chip which is attached to the interconnect substrate and electrically connected to the first interconnect pad array. In the present chip module, the integrated circuit chip is preferably incorporated into a packaged chip having a plurality of chip leads protruding therefrom. The chip leads are electrically connected to respective ones of a plurality of conductive lead pads which are disposed on at least one of the opposed surfaces of the interconnect substrate and are electrically connected to the first interconnect pad array. It is contemplated that the interconnect substrate of the chip array may include a plurality of conductive lead pads disposed on respective ones of the opposed surfaces thereof, with the chip leads of at least two packaged chips being electrically connected to the conductive lead pads on respective ones of the opposed surfaces of the interconnect substrate. As indicated above, the chip array may comprise a memory array or a logic array, with each of the packaged chips preferably being a memory or logic device selected from the group consisting of a TSOP I package, a TSOP II package, a QFP package, and a CSP package or the like.\nThough the interconnect substrate of the chip array of the present chip module may be of a rigid construction, the same is preferably flexible and formed to have a thickness not exceeding about 0.015 inches. Importantly, this minimal thickness of the interconnect substrate, which results in its flexibility, provides the assembled chip array with a thin profile as well. The interconnect substrate is preferably fabricated from a polyamide, though alternative materials possessing similar characteristics may also be employed in relation thereto. Additionally, the interconnect substrate is preferably formed to include a pair of registry notches disposed within respective ones of the peripheral edge segments thereof which extend in generally perpendicular relation to the peripheral edge segment along which the first and second sets of conductive interconnect pads of the first interconnect pad array are extended.\nIn addition to the chip array, the chip module of the present invention comprises a package including a generally square main body and a lid attachable to the main body. The main body itself defines a cavity sized and configured to receive the chip array, and has a generally planar interconnect shelf which extends within the cavity and includes a second interconnect pad array disposed thereon. The second interconnect pad array comprises a third set of conductive interconnect pads which are arranged upon the shelf in a pattern identical to those of the first and second sets of conductive interconnect pads of the first interconnect pad array of the chip array. In the preferred chip module, a total of seventy-one (71) conductive interconnect pads are included in each of the first, second and third sets thereof. The main body of the package also includes a plurality of package leads which protrude therefrom and are electrically connected to the second interconnect pad array. More particularly, the main body preferably includes sixty-eight (68) package leads which extend about the periphery thereof. The package (i.e., the main body and lid) is preferably fabricated from co-fired alumina, though alternative materials possessing similar characteristics may also be employed in relation thereto.\nIn the present chip module, the chip array is insertable into the cavity of the main body such that the first and second interconnect pad arrays are in aligned contact with each other. Subsequent to the insertion of the chip array into the cavity and the electrical connection of the first and second interconnect pad arrays to each other in a prescribed manner, the attachment of the lid to the main body encloses and seals the chip array within the package. In the present chip module, the main body of the package is preferably formed to include an opposed pair of tabs which are insertable into respective ones of the registry notches of the interconnect substrate when the chip array is inserted into the cavity. Advantageously, the receipt of the tabs into the notches facilitates a proper registry between the first and second interconnect pad arrays.\nFurther in accordance with the present invention, there is provided a method of assembling a chip module. The preferred method comprises the initial step of assembling a chip array of the chip module. Such assembly is accomplished by attaching at least one integrated circuit chip, and more particularly at least two packaged chips, to respective ones of the opposed surfaces of the interconnect substrate such that the packaged chips are electrically connected to the first interconnect pad array of the interconnect substrate. Thereafter, the assembled chip array is inspected and tested.\nSubsequent to being tested, the completed chip array is mounted to a package of the chip module. More particularly, the chip array is inserted into the cavity of the main body of the package such that the first interconnect pad array is in aligned contact with the second interconnect pad array provided on the interconnect shelf extending within the cavity of the main body. The first and second interconnect pad arrays are preferably attached to each other via soldering. In addition to the first and second interconnect pad arrays being soldered to each other, the interconnect substrate of the chip array may optionally be bonded to the main body of the package. After the chip array has been properly mounted to the main body of the package, the assembled chip array and main body are inspected and tested. Such inspection includes screening the assembly for proper burn-in, centrifuge, solderablility, etc.\nUpon the completion of its testing, the assembled chip array and main body are preferably vacuum baked to remove moisture from the chip array (i.e., the interconnect substrate and packaged chips). Thereafter, the lid of the package is attached to the main body to enclose and seal the chip array within the package. The attachment of the lid to the main body may be accomplished through the use of a low temperature lid seal or through the use of a low temperature welding process which creates a seam seal between the lid and the main body. Thereafter, the chip module is tested for fine and gross leaks of the enclosed cavity of the package, with the completed chip module then being given a final test."}
{"text": "The present invention relates to devices for extracting particulates from a gaseous stream. More particularly, it relates to an extraction device for extracting liquid particulates from a gaseous stream."}
{"text": "According to a known construction, an electronic device comprises a stacked assembly comprising a base wafer, an integrated circuit chip, mounted onto a front face of this base wafer and including a sensor within its front face, and a protection wafer mounted on top of this front face. Generally, this stacked assembly is surrounded by an encapsulation block.\nSuch an electronic device, generally resulting from a collective fabrication, notably exhibits difficulties in fabrication so that the front face of the encapsulation block runs substantially in the plane of the front face of the protection wafer.\nThere is a need in the art to solve such difficulties."}
{"text": "1. Field of the Invention\nThis invention relates to integrated circuit manufacturing generally and, more particularly, to a method of planarizing a semiconductor wafer prior to the formation of circuits thereon.\n2. Description of the Prior Art\nIn dielectric isolated (DI) wafer fabrication, bonded wafers are becoming less expensive than wafers made by existing grow-handle-and-grind methods of DI wafer fabrication. However, bonded wafers in general may not be sufficiently fiat for advanced integrated circuit manufacturing, e.g., submicron CMOS technology. There exists a need for an inexpensive process for planarizing bonded wafers, preferably using existing grinding/polishing equipment found in many chip manufacturing plants."}
{"text": "The present invention is directed to connecting lengths of well coiled tubing or connecting tools to coiled tubing.\nCoiled tubing is a continuous conduit without joints used in an oil/gas well carried on a spool or reel at the well surface. The coiled tubing can be lowered into an oil and/or gas well and is of a small diameter, for example, one to two inches. Coiled tubing can be used for installing and operating a variety of well accessories, such as disclosed in U.S. Pat. No. 4,844,166. Interconnection of the coiled tubing with other links of coiled tubing and well tools has been accomplished in the past, as described in U.S. Pat. No. 4,682,657, by providing a connector gripping the exterior of the coiled tubing. However, such connectors inherently provide an outside diameter which is greater than the outside diameter of the coiled tubing. This provides difficulty in allowing the coiled tubing and connector to sealingly and freely pass through the injector chains, around the gooseneck, and onto the coiled tubing spool or reel. Furthermore, such connectors are not spoolable. That is, they cannot be wound up on the spool or coiled tubing reel a number of times without losing pressure and tensile integrity.\nThe present invention is directed to a method and apparatus which internally grips and connects to the coiled tubing by an internal ratchettype spoolable connector which avoids the disadvantages of the prior art connections by being spoolable, will sealably pass through packoffs, by having an outer body size no greater than the outside diameter of the coiled tubing, and can withstand a number of cycles of being spooled and unspooled without losing pressure or tensile integrity."}
{"text": "A frequent goal in the design of photographic equipment is a reduction in its overall size and weight. This is particularly true with regard to portable cameras where a reduction in size and weight allows the camera to be more easily transported, as well as more easily manipulated during photographic use.\nA significant factor which dictates the size of photographic equipment is the total conjugate distance, i.e., the distance measured along the optical axis between the object and the image. The total conjugate distance (TC) is a function of the focal length of the lens. In a camera, for example, which houses both an object field such as data card, and a receptor surface, such as photographic film, the ability to reduce the total conjugate distance is important in reducing the overall size of the camera."}
{"text": "Diabetes mellitus is a disorder in which the pancreas cannot create sufficient insulin (Type I or insulin dependent) and/or in which insulin is not effective (Type 2 or non-insulin dependent). In the diabetic state, the victim suffers from high blood sugar, which causes an array of physiological derangements (kidney failure, skin ulcers, or bleeding into the vitreous of the eye) associated with the deterioration of small blood vessels. A hypoglycemic reaction (low blood sugar) may be induced by an inadvertent overdose of insulin, or after a normal dose of insulin or glucose-lowering agent accompanied by extraordinary exercise or insufficient food intake.\nConventionally, a diabetic person carries a self-monitoring blood glucose (SMBG) monitor, which typically requires uncomfortable finger pricking methods. Due to the lack of comfort and convenience, a diabetic will normally only measure his or her glucose level two to four times per day. Unfortunately, these time intervals are spread so far apart that the diabetic will likely find out too late, sometimes incurring dangerous side effects, of a hyperglycemic or hypoglycemic condition. In fact, it is not only unlikely that a diabetic will take a timely SMBG value, but additionally the diabetic will not know if his blood glucose value is going up (higher) or down (lower) based on conventional methods.\nConsequently, a variety of non-invasive, transdermal (e.g., transcutaneous) and/or implantable electrochemical sensors are being developed for continuously detecting and/or quantifying blood glucose values. These devices generally transmit raw or minimally processed data for subsequent analysis at a remote device, which can include a display."}
{"text": "Web site traffic analysis is the process of collecting and analyzing data about a web site's traffic to determine the web site's usage patterns and help measure the site's performance and effectiveness. The analysis may examine the total number of visitors to a site, the portions of a site visited, the duration of the visits, and the details of the interactions with the site, such as search queries generated by the visits. Web site traffic analysis is important to the effective functioning of web sites. By some estimates, the business of web site traffic analysis may gross over a billion dollars a year.\nWeb site traffic analysis may be performed by creating logs with entries for each Hypertext Transfer Protocol (HTTP) request or hit to the web site server. A log entry may be used to generate information about the number of hits, visitors, visitor duration, visitor origin (sub domain, referral link), visitor IP address, browser type and version, platform, cookies, and details of interaction with the web site. Web site traffic analysis may also be performed by creating records of the links or Uniform Resource Locators (URLs) visited at a web site.\nWeb servers may utilize a portal architecture. In this architecture, multiple applications called portlets are integrated into a single, customizable application called a portal. Content provided by multiple portlets may be displayed on a single page, and a user may interact with an individual portlet on a page without necessarily affecting the state of other portlets on the page. The portlets may be local, run by the web server, or remote. A remote portlet for a portal may be operated by a different enterprise than operates the portal. For example, the portlet may provide weather information to a large number of unrelated portals.\nTraditional methods of web site traffic analysis may fail to capture critical information about web site usage of portals. To understand the usage of a portal site, the following data may prove useful for analysis:\nthe name of the portal page being displayed,\nthe names of all the portlets on the page,\nthe states of the portlets on the page (normal, minimized, maximized, etc.),\nmiscellaneous data about state internal to and specific to each portlet on the page, and\nthe user id of the user using the portal.\nThe data listed above may prove unavailable under traditional methods of web site traffic analysis. The web pages generated by portal servers may, for example, not reveal that the pages are portal pages. The pages may consist of Hypertext Markup Language (HTML) pages which contain fragments corresponding to the portlets. The web servers may not place easily extractable information into the page HTML regarding the portal environment, the portal page, or the individual portlets. The portlets themselves may be presented as nested tables with no identification. It may, therefore, prove difficult to determine from the page the details of the interaction between a user and an individual portlet. Further, reliance on records of server URLs may prove problematic. The server URLs may be highly dynamic, reflecting the changing content of portlets in response to user actions. Thus, a logged URL may no longer exist when an analysis is performed. Further, the URL may not expose all relevant information about the state of the page being viewed. Moreover, relying on the server to gather data about web site traffic may miss revisits to cached pages. These pages may be viewed again by retrieval from the cache without making an HTTP request of the server.\nIn addition, under traditional methods of web site traffic analysis, the integration of the collected data may prove difficult. The data about web site traffic may be written to multiple log files without a common identifier to tie them together. Further, the generation of server-created logs by a portal may consume large amounts of the portal's storage or may require intensive computing resources from the portal. Further, a portal may not deliver the logs to the web site traffic analysis provider in a timely fashion."}
{"text": "Mechanical energy comprises a number of forms of energy including, but not limited to kinetic energy. Mechanical energy is manifested in the bodies of humans and animals as a result of their physical processes. Such physical processes include voluntary body movements. Amongst voluntary body movements are gait processes. Gail activities include stepping, walking, running, climbing, jumping, and similar activities. Other voluntary body movements include grasping, reaching, shaking, swinging, stretching, etc. All voluntary body movements are manifested as motion of body members having mass so that all voluntary motor activities develop kinetic energy. Further, voluntary motor activities may impart kinetic energy to peripheral masses engaged with a moving body.\nIt is sometimes desirable to convert mechanical energy to electrical energy. An example is the conversion of kinetic energy into electrical energy as the kinetic energy of a mass moves a magnetic field relative to a conductive coil thereby converting the kinetic energy of the mass to electrical energy by action of electromagnetic induction.\nDevices to convert the kinetic energy manifested in the bodies and peripheral masses engaged with the bodies of humans as a result of their physical processes into electrical energy are not well-developed. Accordingly, it is desirable to provide a device to harvest kinetic energy imparted by voluntary motor activities and convert the harvested mechanical energy into electrical energy."}
{"text": "This invention relates to the diagnosis of Alzheimer Disease in human patients. The invention further relates to assessing the predisposition of a presymptomatic individual to develop the disease. More particularly, the present invention concerns the use of nucleic acid probes to detect gene alterations relating to Alzheimer Disease.\nAlzheimer Disease is a progressive, degenerative disease of the central nervous system, characterized by specific neuropathological lesions. Many neuronal perikarya contain accumulations of neurofibrillary tangles (NFT), which are composed of paired helical filaments (PHF). Neurofibrillary tangles contain a protein structurally different from the amyloid found in extracellular aggregates in amyloid plaque cores and in cerebrovascular lesions of congophilic angiopathy. Similar amyloid lesions occur in such pathological conditions as Down Syndrome, Gerstmann-Straeussler-Scheinker Syndrome, Kuru, and in a minority of patients with Creutzfeldt-Jakob Disease.\nA genetic component is indicated by large scale studies of pedigrees in which ten or more members of four generations have developed a dementia of the Alzheimer type before they were 70 years old. In such cases, it appears that an aberrant gene is transmitted as autosomal dominant. Most cases of Alzheimer Disease occur in people over 80 years of age, and it has been calculated that up to 30% of the Caucasian population over the age of 85 has the disease.\nIt is estimated that Alzheimer Disease currently affects over 2 million elderly people in the U.S. Because the disorder is usually late onset, the number of affected individuals will continue to grow as the elderly population increases in size. There are currently more than 25 million people in the U.S. over the age of 65, and this number will more than double by the year 2025.\nOne of the major biochemical symptoms in the cerebral cortex of Alzheimer patients is the appearance in plaque and blood vessel deposits of a protein known as the amyloid plaque core (APC) protein. This protein is an insoluble, highly aggregating, small polypeptide of relative molecular mass 4,500. APC has been understood to have 40-42 amino acids, and is currently believed more likely to be 40 amino acids in length. This polypeptide is also deposited in the brains of relatively older individuals with trisomy 21 (Down\nSyndrome). Kang et al, Nature 325, 733-736, 1987, propose that the amyloid protein is of neuronal origin and is part of a larger precursor protein. To identify this precusor, a full-length complementary DNA clone coding for the APC protein was isolated and sequenced. The predicted precursor consists of 695 residues and contains features characteristic of glycosylated cell-surface receptors. This cDNA sequence, together with the localization of its gene in chromosome 21, suggests that the cerebral amyloid deposited in Alzheimer Disease and aged Down Syndrome is caused either by aberrant catabolism of a cell-surface receptor, or by aberrant synthesis or processing of its mRNA. In order to determine patterns of expression of the gene in brain, APC cDNA was hybridized to tissue sections of the mouse nervous system. Amyloid precursor protein mRNA was detected in most neurons of the central and peripheral nervous system.\nA means for early and accurate diagnosis of Alzheimer Disease will have a major impact on the progress of research on dementia. Because no peripheral biochemical marker for Alzheimer Disease has been found, a definitive diagnosis of the disorder can be made only if histologic confirmation is obtained by performance of a cerebral biopsy or an autopsy. Despite the fact that Alzheimer Disease is essentially an untreatable condition, there would be significant value to know whether or not an individual with Alzheimer symptoms in fact has the disease. A number of treatable conditions are characterized by symptoms similar to those found in patients with Alzheimer Disease. Some such treatable conditions are brain tumors, thyroid and other endocrine dysfunctions, depression, infection, vitamin and mineral deficiencies, metabolic disorders, unrecognized injuries, and side effects of medication. The medical literature advises physicians to conduct a thorough examination, medical history, and testing in order to rule out reversible disorders before diagnosing the untreatable Alzheimer condition."}
{"text": "FIG. 11 schematically shows a conventional method for manufacturing an optical display device (liquid crystal display). First, an optical member manufacturer produces a long sheet material, which includes an optical member, in the form of a roll (#1), and it is slit into a predetermined size (#2). Examples of the optical member include a polarizing plate for use in a liquid crystal display, a laminated film including a polarizing plate and a retardation plate, and so on. The slit piece of the long material is then cut into a specific length according to the size of an optical display unit (liquid crystal cell) (#3), and the specific-length piece of the sheet material is subjected to an appearance inspection (#4). The finished product is then inspected (#5). Subsequently, the four end faces of the piece of the sheet material are worked (#6). The piece of the sheet material is then subjected to clean packaging in a clean room environment (#7), and packaging for transportation is performed (#8). The piece of the sheet material manufactured as described above is transported to a panel processing manufacturer.\nThe panel processing manufacturer unpacks the piece of the sheet material (#11). An appearance inspection is then performed to check whether scratches, stains or other defects are produced in transit or during unpacking (#12). The piece of the sheet material determined as non-defective in the inspection is then transferred to the next step. An optical display unit (liquid crystal cell) to which the piece of the sheet material will be bonded is previously manufactured and cleaned before the bonding step (#13).\nThe piece of the sheet material and the optical display unit (liquid crystal cell) are bonded together to form an optical display device (#14). The release film is peeled off from the piece of the sheet material so that the pressure-sensitive adhesive can be left, and one side of the optical display unit (liquid crystal cell) is bonded to the surface of the pressure-sensitive adhesive. The other side of the optical display unit (liquid crystal cell) may also be bonded in a similar manner.\nJapanese Patent Application Laid-Open (JP-A) No. 2007-140046 (Patent Literature 1) discloses a manufacturing method, which includes: feeding a sheet material, which includes an optical member, from a roll thereof; detecting whether the sheet material has any defect; cutting the sheet material into pieces based on the result of the detection; subsequently peeling off the release film; and then bonding the cut piece of the optical member to a liquid crystal cell.\nJP-A No. 2005-37416 (Patent Literature 2) discloses a manufacturing method, which includes cutting the members (e.g., a polarizing plate) of a sheet material other than a release film so that the sheet material can be kept continuous by means of the release film and bonding the cut piece of the sheet material to an optical display unit (liquid crystal cell) with the pressure-sensitive adhesive interposed therebetween, while peeling off the release film.    Patent Literature 1: Japanese Patent Application Laid-Open (JP-A) No. 2007-140046.    Patent Literature 2: JP-A No. 2005-37416."}
{"text": "Connecting devices between an upper board and a lower board are known, to achieve a partition modular wall to separate one area from another.\nEach board is therefore placed on a lower board, and sometimes it is fixed to the lower board. Each board bears the weight of all the boards above.\nSometimes, one or more extruded profiles fixed by screws to the upper board and to the lower board are used, further comprising a groove adapted to allow the hooking of objects projecting from the partition wall.\nSuch known walls are accessible from both sides and therefore, for aesthetic reasons, the connection profiles must be as little as possible visible from both sides of the wall.\nAmong other things, the static load on each board is given by the weight of the upper boards and of the other components supported resting thereon.\nThese structures are thus designed to support compressive loads.\nA known device of this type for connecting an upper board to a lower board of a wall resting on a floor involves the use of two separate profiled elements coupled together, fixed to two facing sides of a lower board and an adjacent upper board, respectively. These two elements may include complementary surfaces to align such boards when they are resting on one another.\nSometimes, however, the need is felt to have a fitted panel to be applied hung to a wall, able to allow the engagement to such a panel of cantilever shelves, or other hanging components, positionable in different positions along a horizontal line, or to be able to slide such components horizontally along the wall.\nThe known structures do not allow being hung to an existing wall as they would tend to break down into individual boards since the profiles used are able to resist compressive and not tensile stresses.\nFor example, if an upper board of a known panel as described above were fixed to the wall, all the other lower boards would be hanging on such an upper board and the weight of the same and of the wall units hanging on the panel would weigh thereon. This would cause the separation of the boards of the above described known panels, since the profiles typically used would not withstand such a high tensile stress.\nMoreover, the known structures cannot be hung to an existing wall since they would not be able to withstand the applied loads without deformation on the joints between the boards. In fact, such structures, when hung, would be subjected to a stress distribution different from when they are resting on a floor, and would inevitably tend to assume a deformed configuration not suited to their function.\nMore in detail, it may happen that the application of cantilever loads generates forces torques with respect to the fixing points of such structures to a wall, which inevitably would cause the deformation of the structure itself at the joints between the boards.\nTherefore, the need is felt to provide an extruded connection profile for connecting to each other an upper portion and a lower portion of a fitted panel having hooking slots for support brackets, able to allow the hooking of such a panel to a wall in a hanging manner while preventing such a panel from getting deformed under the action of forces applied to the same at the slots, and preventing such a panel from being broken down due to its own weight and to the weight applied by hanging cupboards suspended on the panel.\nThe need is also felt to provide such a fitted panel adapted to be hung to a wall preventing the deformation under the action of loads."}
{"text": "This section is intended to provide background information to facilitate a better understanding of various technologies described herein. As the section's title implies, this is a discussion of related art. That such art is related in no way implies that it is prior art. The related art may or may not be prior art. It should therefore be understood that the statements in this section are to be read in this light, and not as admissions of prior art.\nA radar may provide marine vessels with the capability to safely navigate a course through unknown bodies of water. As such, a radar may prove to be a useful tool in navigation by identifying neighboring vessels and obstacles in close proximity to a marine vessel."}
{"text": "In the course of the production or refurbishment of switch mechanisms, the forces generated by the switch mechanism must be definitively set and tested. One of the aims is that the forces generated by the switch mechanism do not go above or below certain upper and lower limits. For this purpose, opposing forces that correspond to the forces of the switch blade during later operation are applied to the switch mechanism acting as a test piece while it is performing actuating operations on the test rig.\nThe mechanical dimensions of the switch mechanisms differ very widely according to the area of use, manufacturer and development history. Externally, common to all switch mechanisms is the so-called actuating slider, which transfers the linear movement and force of the mechanism to the switch blade, and (optionally) so-called test sliders, which are likewise connected to the switch blade and make a mechanical indication of the position of the blade to the mechanism possible. For the testing, the switch mechanism is ideally fastened in the same way as it is later on the track (usually by way of a number of screw connections). On account of the different outer dimensions of the mechanisms, the two movable components are located at different positions in relation to the fastening points of the switch mechanism.\nThe document “PT 10K Multi: Variabler Prüftisch für Weichenstellantriebe” [PT 10K Multi: Variable test bench for switch actuating mechanisms], available on the Internet on Mar. 31, 2014 at http://www.probitron.de/fileadmin/pdf/datenblatt_multi.pdf, discloses a variable test bench for switch mechanisms that consists of a basic module and exchangeable test modules. The basic module has a hydraulic system as a force generator, a controller and also measuring and safety equipment. Available for different models of switch mechanisms are individual test modules, which according to the illustration respectively consist of a bench with rollers and position holders that are individually designed for the switch mechanism. The adaptation of the coupling point between the test piece and the test-rig force generator consequently takes place by mounting the test piece on a test module that is individually designed for the test piece.\nA test slider mechanically indicates the actual position of the switch blade to the switch mechanism. For testing the test slider, it would have to be moved at the same time as the actuating slider. However, in this case the movement of the actuating slider and the test slider do not take place synchronously; the length of the displacement is also different. Therefore, the known art dispenses with complete testing of the switch mechanism including the test sliders, since the switch mechanism is tested without the test sliders fitted. The test sliders are only fitted after the test."}
{"text": "1. Field of the Invention\nThe present invention relates to a differential apparatus for vehicles.\n2. Description of Relevant Art\nThere has been disclosed in Japanese Patent Application Laid-Open Publication No. 7-156681 a differential apparatus of an arrangement shown in FIG. 1, where it is designated by reference character 201.\nThe differential apparatus 201 includes a differential mechanism 207 for distributing drive power from an engine to a pair of left and right drive axles 203, 205, a pair of left and right drive power control mechanisms 215 each comprising a combination of a multi-plate clutch 211 and a speed change mechanism 213 for a serial connection between a differential casing 209 of the mechanism 207 and the left or right axle 203 or 205, and a hydraulic actuator 220 for actuating the clutch 211 to be internally coupled for the connection, and a hydraulic controller 230 for controlling the actuator 220 to have the clutch 211 coupled.\nIf the multi-plate clutch 211 of the left drive power control mechanism 215 is coupled, drive power distributed to the left axle 203 has its fraction fed back via the speed change mechanism 213 and that clutch 211 to the differential casing 209, with the more distributed drive power at the right axle 205 end. To the contrary, if the multi-plate clutch 211 of the right drive power control mechanism 215 is coupled, drive power distributed to the right axle 205 has its fraction fed back to via the speed change mechanism 213 and this clutch 211 to the differential casing 209, with the more distributed drive power at the left axle 203 end.\nThe controller 230 is thus adaptive for coupling the multi-plate clutch 211 of the left or right drive power control mechanism 215 to control a yaw momentum of a vehicle body, securing a stable straight forward travel as well as a comfortable steering to turn.\nMore specifically, when the vehicle body in travel has entered a yawing due to a slipping road wheel at the left or right axle 203 or 205 end, the drive power control mechanisms 215 are controlled to stop the slipping.\nIt therefore is impossible for the differential apparatus 201 to limit a developing difference at the differential mechanism 207 before the road wheel slips, and to effect a prevention of a wheel slippage before a development of yaw momentum on the vehicle body in a straight forward travel.\nThe respective multi-plate clutches 211, which are arranged between the differential casing 209 and the left and right axle 203, 205, might have been operated for a use of their frictional torque to limit a difference development at the differential mechanism 207. However, if the clutches 211 were both coupled for the difference limitation, those fractions of drive power fed back from the left and right axle 203, 205 ends would have canceled each other. Still worse, during such the limitation, the multi-plate clutches 211 should have stood with a big burden due to slippage, resulting in a reduced service life.\nIn addition, associated drive torque to be fed back is subjected to an amplification at the speed change mechanism 213 adversely affecting the vehicle body in performance, rendering control of difference limitation very difficult.\nThe differential apparatus 201 may involve a difference limiting mechanism, with an arrangement the more complicated, large sized, increased in weight and dear of cost.\nThere has been disclosed in Japanese Patent Application Laid-Open Publication No. 5-345535 another differential apparatus shown in FIG. 2, where it is designated by reference character 301.\nThe differential apparatus 301 comprises an input part 302 for inputting drive power from an engine, a differential mechanism 303 of a planetary gear type, and a drive power control mechanism 305.\nThe drive power control mechanism 305 includes a speed changing mechanism 313 as an integrated combination of a speed increasing mechanism 313a for transmitting a fraction of drive power of a left drive axle 307 to a first intermediate shaft 311, with an increased revolution speed, and a speed decreasing mechanism 313b for transmitting a fraction of drive power of the left axle 307 to a second intermediate shaft 315, with a decreased revolution speed. The control mechanism 305 further includes a hydraulically actuatable first multi-plate clutch 317 to be coupled for interconnection between the first intermediate shaft 311 and a right drive axle 309, and a hydraulically actuatable second multi-plate clutch 319 to be coupled for interconnection between the second intermediate shaft 315 and the right axle 309.\nIn other words, the drive power control mechanism 305 is adaptive for a serial connection between the left and right axles 307, 309 with the first and second multi-plate clutches 317, 319 controlled to be either coupled by a hydraulic controller 330 to achieve a secured stable straight forward travel as well as a comfortable steering to turn.\nThis differential apparatus 301 has similar issues to that differential apparatus 201."}
{"text": "The present invention relates to an improved bee attracting composition, an apparatus to attract bees formed from that bee attracting composition, and a method and apparatus using Applicant\"\"s improved bee attracting composition to attract and immobilize bees.\nThe immigration of Africanized bees, sometimes referred to as xe2x80x9ckiller bees,xe2x80x9d into the United States, and widespread fear of these more defensive bees has increased the public\"\"s demand for removal of swarms and feral colonies of bees. At the same time it has made bee removals more problematic. Compared to European bees, Africanized bee colonies swarm more frequently and occur at higher density, so that swarms are more frequently encountered. When swarms or colonies are discovered, there is an increased likelihood that the bees will sting, and a much increased perception that they are dangerous. At the same time, there are fewer beekeepers willing or permitted to perform bee removals. Pest control specialists, and in some areas public agencies, are removing increasing numbers of bee colonies.\nSwarms may be removed by hiving them or by spraying them with insecticides, including soap or detergent solutions. A major problem with either technique is that xe2x80x9clostxe2x80x9d bees remain at large after most of the bees have been killed or hived. Most of these are bees that were scouting nest sites or foraging away from the swarm cluster when control efforts began. It has been estimated that about 5% of the bees in a swarm participated in scouting one nest site, so in a swarm of 10,000 or more bees, the scouts can comprise several hundred bees. xe2x80x9cLostxe2x80x9d bees from a removed swarm remain near the swarm site, but fly around a great deal searching for their queen. These bees can survive for at least several days, and are likely to run short of food, which may make them more likely to sting.\nMuch has been learned about the behaviors of insects, including bees, in recent years. It is now understood that much of bee behavior arises from use of chemical substances, and mixtures of chemical substances, called xe2x80x9cpheromones.xe2x80x9d A pheromone is a substance secreted by an animal that causes a specific reaction by another individual of the same species. Many bee activities can be explained as the effect of various pheromones.\nApplicants\"\"novel invention includes a bee attracting composition formed by combining aliphatic and aromatic organic compounds which simulate various naturally-occuring honey bee pheromones. Applicant\"\"s bee attracting composition can be formed into a solid apparatus which can be used to attract bees.\nApplicant\"\"s invention further includes an apparatus and method to attract and immobilize bees. This apparatus is formed from a substrate, an adhesive composition, and Applicant\"\"s bee attracting composition. The treated substrate can be folded to form a synthetic bee hive. Bees are lured in the interior of the synthetic hive, are drawn to the adhesive component by the bee attracting composition, and become permanently immobilized upon contact with the adhesive."}
{"text": "Most ion guides consist of multipole structures, extending longitudinally, having rod-shaped pole pieces. Their disadvantage is that they do not actively drive the ions forward without complicated additional measures. For this reason, ion guides consisting of stacked diaphragms with circular apertures (known as “stacked rings”) are sometimes used for special purposes; an axial DC or modulated potential gradient permits the ions to be driven forward actively. Examples of this include ion funnels used to capture the ions from a gas flowing into the vacuum, collision cells with diaphragms of constant internal diameter and with active forward drive (“ion tunnels”), and ion packeting equipment using a traveling wave field to provide a forward drive of an ion beam with a desired time profile of ion density.\nFor example, U.S. Pat. No. 6,107,628 (R. D. Smith and S. A. Shaffer) elucidates an arrangement of an ion funnel in which ions are extracted from a stream of gas and are directed to the output opening that leads to the next differential pump stage. The ion yield is significantly greater than when simple skimmer diaphragms are used. This ion funnel represents a special case of the general description of ion guides in U.S. Pat. No. 5,572,035 (J. Franzen), in which, among other embodiments, arrangements of stacked rings have already been described, operated at radio frequency voltages (RF) with an axial direct current (DC) potential gradient and with both cylindrical or conical internal open space.\nDiaphragm stacks in the form of ion funnels are being used more and more frequently instead of the gas skimmers usually applied. The ion funnel consists of a package of coaxially arranged diaphragms having circular apertures, in which the diameters of the circular holes diminishes increasingly toward the central exit hole that leads into the next chamber. The diaphragms are stacked with relatively small spaces between the apertured diaphragms. This creates the shape of a funnel inside the stack of diaphragms. Gas with entrained ions from an ion source that is external to the vacuum is blown through an inlet opening into the vacuum system, or through an inlet capillary, into the open ion funnel. The wall of the ion funnel is highly permeable to gas, as it is formed from the faces of the apertured diaphragms with the open spaces between them. The gas escapes through the spaces between the apertured diaphragms, and is removed by a vacuum pump. Only very little gas enters the next chamber of the differential pump system through the small exit opening.\nBoth phases of an RF voltage (several hundred kilohertz up to several megahertz; a few hundred volts) are applied alternately to the apertured diaphragms. This repels the ions from the inner funnel wall. The method of operation and the effect of this repelling pseudopotential are described in detail in the quoted patent specification, U.S. Pat. No. 5,572,035. The ions are thus prevented from being drawn away through the intermediate spaces between the apertured diaphragms by the escaping gas stream. The ions are separated out. In addition, a graded DC voltage (a few tens of volts in total) is applied to the apertured diaphragms to create a potential gradient along the axis of the stack of diaphragms. This forces the mobile ions through the highly rarefied gas in the ion funnel toward the exit hole.\nThe system of annular diaphragms, including the ion funnel, has the advantage of actively driving the ions forward to the exit of the annular diaphragm system. They have, however, the disadvantage that, even in the presence of a cooling damping gas, the ions are not collected along the axis of the annular diaphragm system, since the pseudo-force that repels the ions only exerts its effect close to the outer wall of the cylinder or cone created by the diaphragm openings. The ions therefore fill the entire internal space of the cylinder or cone. If this space is densely filled with ions, the Coulomb repulsion (“space charge effect”) will even increasingly drive the ions against the pseudopotential wall, whereas the part of the internal space that is close to the axis has a lower ion density.\nThe repelling effect of the walls around the interior space is, moreover, different for ions with different specific masses. “Specific mass” here refers to the ratio of mass to charge. For heavy ions (ions with high specific masses) the ions are only reflected when close to the wall, whereas lighter ions are reflected at a greater distance from the wall.\nThe embodiment of the ion funnel that has become familiar is particularly disadvantageous from this point of view. The published embodiment has the disadvantage that only a relatively narrow band of specific masses passes through. If the diaphragm openings at the final exit are very small, the pseudopotential from the walls of the narrow channel overlap, and the rise in the overlapping pseudopotential reflects light ions back into the funnel; they cannot leave the funnel. Additionally, the pseudopotential along the axis of the ion funnel displays ripples with potential wells, which collect ions inside and can only be emptied if the axial potential gradient has a certain minimum value. On the other hand, too much gas will enter into the next differential pump stage if the diaphragm openings at the final exit of the funnel are too large. If large diaphragm openings are followed by an extraction lens with narrow openings to extract the ions from the ion funnel, the heavy ions cannot be extracted if the space charge is high, since they will be driven outward to the walls of the funnel and escape the drawing field of the extraction lens, which can only effectively extract ions out from the axis.\nOne solution that is already known is an ion funnel that consists of annular diaphragms each of which is divided into four quadrants, wherein the four quadrants of each annular diaphragm alternately carry the two phases of the RF voltage. The next annular diaphragm then carries the phases of the RF voltage crosswise. Manufacture of this quadrant funnel is, however, extraordinarily difficult and expensive.\nThe individual guiding elements of the mass spectrometer through which the ions are to pass generally have very sharply defined acceptance cross sections for incoming ions, different as well for the distribution of directions as for the energy distribution of the ions available in the ion beam. The beam cross section in particular can generate high or low transmission of ions into the next section. For example, the literature describes a very narrow, elliptical acceptance cross section for a quadrupole filter. The narrow acceptance cross section extends between the two pole pieces that carry the DC voltage that attracts the ions. In contrast, a time-of-flight mass spectrometer with orthogonal injection of the ions into an ion pulser requires a very narrow ion beam, close to the axis, and with the most homogeneous distribution of directions and energies that can be achieved. These requirements cannot be satisfied by the stack of annular diaphragms constructed in the manner known at present, even though the possibility of actively driving the ions in the axial direction is a strong factor in favor of the use of diaphragm stacks."}
{"text": "The present invention relates to entry rug cleaning assemblies which remove loose soil from the soles of the shoes of persons passing thereover and in which the surface of the rug is cleaned periodically.\nIn applicant's prior U.S. Pat. No. 3,526,015 granted Sept. 1, 1970, there is disclosed an entry rug cleaning mechanism which has enjoyed commercial acceptance and which has proven quite useful in cleaning and disinfecting the soles of shoes of persons passing thereover into a building. In this installation, the rug is mounted so that it may be rotated about a closed loop, and within an adjacent housing disposed below the traffic surface are vacuum cleaning and washing installations for cleaning the surface of the rug passing thereover.\nSelf cleaning entry rug installations of this type are extremely desirable from the standpoint that fresh rug surface may be provided to effectively remove soil from the soles of the persons passing thereover and they are able to provide a relatively clean rug surface from an aesthetic standpoint as compared with fixed rug installations for this purpose.\nAlthough the installations of applicant's prior patent have provided substantial benefits, location of the cleaning components outside the belt loop increases the space required for the installation and does not provide optimum cleaning action since the area available for cleaning action is always a vertical portion of the belt loop. Moreover, frequent servicing of the cleaning components has been required.\nIn addition, replacement of the rug was frequently difficult and time consuming, and the tendency of the rug to stretch required constant adjustment of the tension.\nAccordingly, it is an object of the present invention to provide a novel self-cleaning entry rug assembly providing enhanced cleaning action and adapted for facile servicing and relatively trouble free operation.\nIt is also an object to provide such an entry rug assembly in which the components may be readily fabricated and assembled from relatively durable elements to provide a rugged installation.\nAnother object is to provide such an entry rug assembly which is adapted to fully automatic control, semi-automatic control, or manual operation.\nStill another object is to provide such an assembly which may be linked to presence sensors or approach sensors to limit operation during the time when persons are passing thereover."}
{"text": "The present invention refers to a method of sterilizing a package material. More precisely the invention refers to a method of sterilizing a package material having a plastic surface by means of hydrogen peroxide in gaseous phase being supplied to the previously heated plastic surface.\nSterile containers which are ready to be filled can be obtained in various ways. One method of sterilization implies that the plastic containers are sealed directly in connection with the manufacturing procedure, e.g. when blow moulded. Due to the collapse of the container at a subsequent lowering of the temperature the sealing in this case results in considerable problems. The container must therefore be stored at an elevated temperature since the volume of air enclosed in the container will decrease if the temperature is decreased. Such a storage is both cumbersome and costly. In order to avoid this, among other materials hydrogen peroxide is nowadays utilized in a liquid or gaseous phase, the inside of the containers then being sterilized.\nHydrogen peroxide is a very potent but not especially stable oxidant, and an advantage is that only plain water is formed as by-product. However, the half-life of hydrogen peroxide is very long, and the sterilizing agent has thus to be removed after it has been allowed to exert its effect.\nHydrogen peroxide is supposed to exert its antimicrobial activity through its powerful oxidizing effect by acting on components in the cytoplasmic membrane of the microorganisms. In the sterilizing procedure hydrogen peroxide is used as a water solution which is sprayed into a container. When the container is subsequently heated the hydrogen peroxide is gasified and can be ventilated away after having exerted its effect.\nAccording to an alternative sterilizing procedure a container is in a first step subjected to hydrogen peroxide in gaseous phase, which is allowed to condensate on the inside of the container. By this procedure a homogeneous coating of hydrogen peroxide is obtained, which completely covers the package material. In a subsequent step the condensated gas is again gasified and is finally ventilated away. By this procedure an acceptable death rate of microorganisms for commercial sterilization can be achieved.\nWhen for example plastic containers are sterilized these are in practice coated with a commercial standard solution, called xe2x80x9cOxonia Aktivxe2x80x9d, which among other things contains hydrogen peroxide and peracetic acid. The solution is allowed to act for some minutes, and the containers are then rinsed with clean sterile water. Such a procedure is both complicated and time consuming.\nA further problem is that with certain plastics in the package material, above all polyester, residues of hydrogen peroxide are obtained which are adsorbed on the inside of the container as a result of spraying, coating or condensation. They can afterwards have a negative influence on the filling material, for example by changing the colour. The authorities have established limit values which must not be exceeded for residues of hydrogen peroxide.\nIt is rather unusual to treat such a package material as a polyester aseptically since this type of material mostly is used for soft drinks and the like. However, there exists a great interest from the market for an aseptic bottle of polyethylene terephtalate (PET).\nIn order to minimize the risk of residual product formation the hydrogen peroxide can be ventilated away, e.g. at the elevated temperature or by sterile air being blown into the container. Such a procedure is costly and it would thus be desirable if shorter ventilation times could be used, which would result in a lower cost per container.\nThe purpose of the invention is to eliminate or reduce the drawbacks mentioned above which arise during a sterilization by means of hydrogen peroxide according to the state of the art.\nIn order to achieve this purpose the method according to the invention has obtained the characterizing features of claim 1.\nA sterilizing process with a condensation of hydrogen peroxide comprises four phases of different time intervals: a heating phase, a gassing phase, a holding phase, and a ventilation phase.\nDuring the heating phase the container is heated to a suitable temperature. The heating phase is not necessary. However, it reduces the amount of condensate which later on has to be ventilated away.\nDuring the gassing phase gaseous hydrogen peroxide is injected into the previously heated container and over the surface which is to be exposed to the treatment, onto which the gas condensates, and a condensate in the form of droplets is formed. The gas used is a mixture of air and a gasified solution of hydrogen peroxide. The solution of hydrogen peroxide is supplied as a commercial grade which consists of about 65% water and about 35% hydrogen peroxide. However, other concentrations can of course be used in dependence of other parameters.\nDuring the holding phase the condensate is allowed to exhibit its effect, and during the ventilation phase the container is ventilated, the condensate again being transformed into gaseous phase and removed from the container. In this connection hot air is used which has been pre-filtered and sterilized by means of filtration.\nExperiments were performed with the purpose of achieving an equally effective sterilization as that which allows the hydrogen peroxide used to exert its effect in the form of a condensated gas. In these experiments the hydrogen peroxide was according to the invention allowed to remain in the gaseous phase during the entire sterilization procedure.\nIn comparison with the sterilization by means of condensation described above the method according to the invention implies briefly that the containers are sterilized with gasified hydrogen peroxide, the containers first being heated to a temperature which is higher than the dew point of the gas used. The containers are then blown to dryness with hot sterile air.\nDuring the sterilization procedure the containers are thus for a certain period of time exposed to gaseous hydrogen peroxide of a specific temperature, dew point and flow rate. Apart from the hydrogen peroxide concentration these parameters are important in order to achieve an effective killing. In this connection the dew point is the temperature the gaseous hydrogen peroxide must be cooled to in order to obtain a condensation, a water solution of hydrogen peroxide being separated. This temperature is dependent on the hydrogen peroxide content in the gas. The dew point is adjusted to be used in an existing sterilizing equipment.\nPreliminary experiments have shown that the amount of residual hydrogen peroxide per unit area varied after contact with different types of package material. Certain problems were obtained with polyester as a package material, and thus the experiments were concentrated on such a material by using bottles of polyethylene terephftalate. However, other types of containers which are ready to be filled were also studied."}
{"text": "Field of the Invention\nThe invention relates to an electronic device having an operating mode and an energy saving standby mode, and to a method for switching between the two modes.\nThe field of the present invention preferably covers receivers and transmitter/receivers (transceivers) that are supplied with voltage by a battery. Devices and methods for energy saving are used to prevent the batteries in these appliances from quickly becoming dead or empty. In order to save electrical energy in states in which there is no need for the electronic device to be operated, this appliance can be switched from a normal operating mode to an energy saving standby mode. In this case, all those parts of the appliance that result in a high current draw are generally switched off. The only parts that remain active are those used for switching back to the operating mode. This switching-back process is generally carried out automatically by identifying a valid external signal. This drastically reduces the overall current draw of the electronic device. Since, in the case of battery-powered electronic devices, the time in the standby mode is generally many times longer than the time in the operating mode, it is desirable for the appliance to have a very sparing current draw in the standby mode, in order to maximize the useful life of the battery.\nAn electronic device of this generic type is known from Issued European Patent EP 0 554 386 B1. The receiver includes an input-side antenna element for receiving a signal carrying information, in this case in the form of an RF (radio frequency) radio signal.\nFor the purposes of the present invention, the term antenna element is defined as an input-side element of the electronic device, which converts radiation energy to energy carried by cables for the receiving unit. In principle, radiowaves, microwaves, lightwaves or waves in some other frequency spectrum that is suitable for transporting information can be received by the antenna element as radiation energy. By way of example, the “antenna element” may thus also be a photoelement in order to receive lightwaves. If energy supplied by cables is already available on the input side, then it may even be possible to dispense entirely with the antenna element, replacing it by a suitable interface—for example by a socket arrangement.\nIn the case of the present receiver, a receiving unit that is connected downstream from the antenna element has a demodulator for recovering the AF (audio frequency) radio signal as basic information that has been modulated on the RF signal in accordance with the modulation method used at the transmitter end. An optical indication for the operating parameters, various control elements for manually adjusting the receiver, as well as a loudspeaker for emitting the AF radio signal are controlled by a central controller, which controls their interaction. A battery is provided as the voltage supply unit for providing the supply voltage required for operating the receiver. In order to achieve a long battery life, the receiver has an operating mode and a standby mode. Switching devices for switching between the two modes are integrated into the central controller.\nSwitching from the operating mode to the standby mode is carried out automatically when a noise switching-off signal is present. In the standby mode, the receiver is not ready to operate. The system is woken up from the standby mode by periodically switching on the operating mode at defined time intervals. A timer, whose time consumption is as low as possible, runs in the standby mode for this purpose. When the receiver is ready to operate, the controller thus very quickly checks whether a valid RF radio signal is being received on the input side. If a noise switching-off signal is present instead of this, then the receiver returns to the standby mode again. This method has the disadvantage that the process of switching on the operating mode on a sample, time-controlled basis, which is often unsuccessful, results in quite a high current draw. Attempts have been made to minimize the current draw by identifying the presence of a signal on the input side as quickly as possible, in order to make it possible to return to the energy saving standby mode as quickly as possible where necessary. With this method, when information is not transmitted continuously, that is to say it is transmitted in a pulsed form, it is still possible for a problem to occur when the transmitter is transmitting during a time in which the receiver is not ready to operate, so that the receiver cannot reliably receive the signal carrying the information. In order to preclude this problem, it is necessary either for the transmitter to transmit for a long enough period, or for the time intervals between the waking-up processes to be matched to the transmitted pulse length.\nA method for waking up from the standby mode is also generally known, in which the problems discussed above are avoided by ensuring that the electronic device is ready to operate all of the time, despite being in a standby mode with restricted functionality. In this case, circuit components whose current draw is high are switched off in the standby mode. The only circuit components which are kept in operation are those which are used for monitoring the status of devices on the input side, such as receivers, sensors, detectors. In this way, it is possible to confirm when there is once again a requirement to switch on the other circuit components whose current draw is high, in order to initiate the process of waking up from the standby mode to the operating mode.\nBy way of example, the receiver sensitivity can be reduced in the standby mode. The receiver can then distinguish only roughly between the system's own signals and external signals, but draws less current in the process. As soon as a signal arrives which overrides the selection to maintain the standby mode, the receiver changes fully or partially to the operating mode in order to check the validity of the signal. In comparison with the method described initially, the transmitter may admittedly transmit signals in a pulse manner in this method; however, since it is necessary for the circuit components to be permanently active, the current draw is generally greater. In this case, owing to the reduced sensitivity, it is also difficult to distinguish between disturbances and the system's own signals, which leads to spurious changes to the operating mode. This results in a disadvantageous increase in the energy consumption."}
{"text": "1. Field of the Invention\nThe present inventions relate to processes comprising crosslinking polyethylene or using crosslinked polyethylene, and to products obtained therewith.\n2. Description of the Related Art\nUltrahigh molecular weight polyethylene (UHMW PE) is well known for, among other properties, its chemical resistance, low friction coefficient, high toughness, and resistance against wear. As a result, it has found numerous applications in demanding environments where some or all of the above properties are required. Equally well known is the intractability of this polymer, which is due to the fact that UHMW PE above its crystalline melting temperature, does not form a fluid phase that is of a viscosity that permits melt-processing techniques used with common thermoplastic polymers.\nA negative consequence of the fact that above its crystalline melting temperature UHMW PE merely passes into a viscoelastic state, is that the original polymer particles soften but tend to substantially retain their shape. Even after relatively long sintering times at relatively high temperatures, much of the initial powder morphology may still remain. The latter, in fact, is thought to contribute to wear and ultimate failure of, for example, artificial implants made from this polymer. Attempts have been made to design processes that reduce wear. These processes are however relatively cumbersome, e.g. are ill defined, require melting the UHMW PE more than once, are time consuming, and/or substantially negatively influence other properties of the UHMWE PE such as, e.g., crack resistance."}
{"text": "1. Field of the Invention\nThis invention relates to silver halide photographic light-sensitive elements and more particularly to a method of forming images wherein the silver halide photographic light-sensitive element contains matting agents.\nFinely divided materials with a mean particle size of from about 1 to about 10 micrometers are commonly used as matting agents to provide a rough surface to photographic elements. See, for example, U.S. Pat. Nos. 4,885,219 and 4,022,622. Further, U.S. Pat. Nos. 4,396,706 and 5,057,407 provide matte particles and techniques in order to increase the adhesion of the particles to the photographic element during processing of the element. In addition to the problems expressed in the previously recited references, printer dusting is also an objectionable problem associated with inadequate matte adhesion.\nIt has been heretofore known to use silica particles having a average particle size of 1 to 10 micrometers as mattes for use in photographic elements. In this regard, note U.S. Pat. Nos. 4,022,622; 2,976,250; 3,920,456; 4,409,322; and 4,396,706. The use of silica particles as matting agents in photographic films suffer from a number of disadvantages: they produce an objectionable, slightly milky appearance, their average grain size cannot be closely controlled, they adhere to wall surfaces and therefore give rise to extensive cleaning of equipment and increased labor costs and they settle out in the coating device and supply pipes, thus rendering impossible the precise metering of given quantities to the coating formulation. In addition to the above, the use of colloidal silica in conjunction with matte is disclosed in U.S. Pat. Nos. 4,975,363; 4,914,012; and 4,232,117.\nU.S. Pat. No. 4,235,959 suggests the use of matte particles prepared by condensing in an aqueous medium urea and formaldehyde while vigorously stirring the mixture until particles comprising urea-formaldehyde resin and silica are formed wherein the silica is embedded within the resin matrix."}
{"text": "Broadband wireless is expected to be one of the main drivers of the telecommunications industry. There is a substantial increase in demand for broadband connectivity, with personal broadband being the key growth engine for mobile wireless broadband networks.\nCommunication in such networks is generally divided between access and backhaul. An access network is the air interface network providing traffic communication between mobile terminals (subscribers) and their associated access points (base stations), while a backhaul network is the air interface network providing traffic communication between the various base stations and a core network. The networks may be arranged to transfer data alone, as in Wi-Fi networks, or may be arranged for triple play services video, audio and data), typically WiMax (or other competitive technology, such as 3GPP-LTE). In conventional systems, the access network and the backhaul network each require their own separate transmission equipment, antennas, etc. at great cost to the operator.\nOne example of a conventional backhaul network is connecting wireless base stations to corresponding core mobile networks (ASN GateWay, AAA servers, etc). The choice of backhaul technology must take into account such parameters as capacity, cost and coverage. Base station backhaul typically is performed via wired infrastructure (e.g., E1/T1 leased lines), or via wireless Point-to-point (PTP) microwave links to each base station, which is expensive to deploy (equipment and installation). In particular, due to the direct, uninterrupted line-of-sight requirements of the wireless backhaul equipment, the backhaul components of conventional base stations require strategic deployment location on high and expensive towers.\nMobile WiMAX, as defined in IEEE Standard 802.16e-2005 Standardization for WiMAX, was originally designed to provide mobile broadband access for mobile devices, i.e., broadband wireless data-optimized technology, providing carrier-grade triple play services using a variety of user devices (such as laptops, PDAs, handheld devices, smart phones, etc.). A complete mobile WiMAX Radio Access Network (RAN) requires deployment of massive infrastructure, including base station sites with high towers, base station equipment, antennas, and a separate backhaul network, as described above.\nThe traditional approach for mobile WiMAX network infrastructure deployment is similar to that of cellular phone networks. The network is based on macro-cell deployment, that is, the base stations, radios and antennas are installed on top of high towers, transmitting at high power, so as to maximize the base station coverage area. In order to optimize the cost, the goal is to minimize the number of sites. This can be achieved by deploying more powerful base station equipment for increasing the cell range (e.g., high power radios, multiple radios on each sector with smart antenna techniques), resulting in more expensive base station equipment. However, for a broadband wireless network deployment, this approach is adequate mainly for the coverage phase, when a relatively small number of subscribers share the cell capacity. As the cell coverage area is large, covering a large number of potential subscribers, additional subscribers in each area can rapidly be blocked due to limited base-station capacity.\nOne proposal for increasing the traffic capacity of the base station, while not greatly increasing the interference caused with neighboring cells, is to use several directional antennae on a base station, each pointing in a different direction. In this way, it is possible to “sectorize” the base station so that several different sectors are served from the same location. (In any given direction, only a small number of frequencies are utilized). However, directional antennas are relatively expensive.\nThere are also known outdoor Wi-Fi networks, deployed mainly according to outdoor Wi-Fi mesh technology. The typical Wi-Fi setup contains one or more Access Points (APs), which is the equivalent terminology to Base Station in WiMax, having relatively limited range, deployed along telephone poles, street poles, electricity poles and rooftops. Due to the access point unit's smaller coverage range, a large number of access point units are required to cover a given area. Conventional outdoor Wi-Fi access point units require costly power amplifiers in each Wi-Fi AP unit to extend the coverage range. In addition, conventional Wi-Fi networks operate only on unlicensed bands and suffer from severe interference and difficult radio-planning issues.\nFurthermore, in the micro/pico-cell deployment approach of conventional Wi-Fi-mesh networks, due to multiple access point nodes in the network, backhauling becomes more complicated and costly. Backhauling each node via wired lines (E1/T1 or DSL) is impractical in a dense deployment of nodes. On the other hand, backhauling each node via traditional wireless PTP microwave links is expensive due to costly equipment and installation costs and not feasible to deploy on telephone poles, street poles, electricity poles, etc. In Wi-Fi, like in WiMAX, PTP microwave links require high towers to achieve a clear line-of-sight between nodes. In addition, when the network load is increased, the backhaul network losses drastically degrade the overall network performance (capacity and latency).\nIn traditional Point-to-Point (PTP) microwave backhaul operating in licensed bands or using unlicensed bands, OFDM (Orthogonal Frequency Division Multiplexing) or single carrier technology (constant power with a fixed modulation scheme) are typically employed. In OFDM, the channel bandwidth is divided into multiple concurrent parallel transmissions on several frequencies. However, during each time slot, there must be transmission over every frequency in the bandwidth. Thus, there is no granulation to permit correction of local interference, and, if there is a problem with transmission on one frequency, the entire transmission can be lost due to lack of frequency diversity, so the Signal to Noise Ratio (SNR) of a link (between two base stations) falls on the entire link.\nConsequently, there is a long felt need for a wireless mobile broadband network that is relatively low in cost and provides an in-band backhaul network having interference mitigation. In particular, it would be useful to have such a network with improved PTP communication having high frequency diversity, variable modulation and coding, MIMO and dynamically adapted Beam Forming, according to the instantaneous link conditions."}
{"text": "The present invention relates to a PON (Passive Optical Network) system in which a plurality of subscribers' devices share optical transmission lines.\nThe PON generally comprises one OLT (Optical Line Terminal) and a plurality of ONUs (Optical Network Units) and transforms signals from terminals such as PCs (Personal Computers) connected to ONUs into optical signals before sending them to the OLT via optical fibers. The optical fibers from the plurality of ONUs are connected together by an optical splitter which optically multiplexes (time-division multiplexes) the optical signals and sends them to the OLT.\nThe length of optical fibers between ONUs and OLT is specified by Chapter 8 and 9 in ITU-T Recommendations G.984.1 to be in the range of, for example, 0-20 km, 20-40 km or 40-60 km, and each ONU is installed at any desired distance from the OLT within the above ranges. Therefore, transmission delays of optical signals between OLT and respective ONUs differ depending on the length of optical fiber. Without considering the transmission delays, there is a possibility of the optical signals output from ONUs colliding or interfering with one another when they are optically multiplexed by the optical splitter.\nTo deal with this problem, the OLT uses a ranging method as defined in Chapter 10 in ITU-T Recommendations G.984.3 to adjust the delays of output signals from ONUs to make the respective ONUs appear to be located at equal distances, for instance 20 km, from the OLT in order to prevent the optical signals from ONUs from interfering with one another. Further, to allocate a communication band of one optical fiber to as many ONUs as possible according to user requests, a DBA (Dynamic Bandwidth Allocation) method whereby the OLT allocates an upstream bandwidth for each ONU (data transmission position/time) is also defined in ITU-T Recommendations G.983.4. A transmission band control based on this method is also being practiced.\nFor example, according to the specifications of Chapter 8.2 of ITU-T Recommendations G.984.3, signals transmitted from a plurality of ONUs to the OLT are called upstream signals and comprised of a preamble, a delimiter and a payload signal. As shown in FIG. 8-2 of Chapter 8 of the same recommendations, a guard time is set immediately before the upstream signal to avoid a possible collision with a preceding burst signal. According to the specifications of Chapter 8.1 of the same recommendations, signals sent from the OLT to the plurality of ONUs are called downstream signals and comprised of a frame synchronization pattern, a PLOAM area, a US Bandwidth MAP area and a frame payload.\nAs shown in Chapter 8.1.3.6 of the same recommendations, an area called a US Bandwidth MAP is used to specify a send permission timing for the upstream signal from each ONU. The US Bandwidth MAP area has a Start value representing a start of the send permission timing and an End value representing an end of that timing, both specified in bytes. These values are also called grant values as they permit the transmission. A difference between the End value and the next Start value is an area where there is no upstream signal, and corresponds to the guard time. Each ONU can be assigned a plurality of bandwidth allocation units called T-CONT and the send permission timing for the upstream signal is specified for each T-CONT."}
{"text": "The invention relates to a dead-front electrical connector which is inexpensive to manufacture and safe to use.\nThere are two broad categories of electrical plug connectors: live-front and dead-front type. In the live-front type, the screws for securing the cord conductors to the terminals are reached through the face or front of the connector. After the screws are tightened the front is usually covered by a removable insulating disc. This type of connector has few parts and is simple and inexpensive. It requires, however, a high level of care in wiring and maintenance to ensure safety. One common danger of such connectors is that the removable insulating disc covering the front may become damaged or lost and the live terminals may thereby be exposed. In connectors of the dead-front type there is a heavy insulating base at the front of the connector and the terminal screws are accessed from behind this base. The wiring area is usually fully enclosed after wiring. Such connectors can accommodate greater variations in the quality of wiring and maintenance without loss of safety or service. However, prior art electrical connectors of the dead-front type generally have had more parts and have been more costly to produce than connectors of the live-front type.\nOne of the many ways of connecting insulated conductors to terminals of wiring devices is the so-called insulation displacement technique which typically involves pushing a conductor into a terminal slot by a pair of hand pliers such that the slot cuts through the insulation and makes electrical contact with the conductor wire. This technique typically requires skillful operators and convenient access to the terminal slots permitting the use of hand pliers.\nThe invented connector is of the dead-front type and makes use of a particularly convenient and effective combination of pivoted side covers, locking means and insulation-displacement terminals. One embodiment of the connector comprises a housing made up of a dead-front base and a pair of side covers having front ends pivotally connected to the base by web hinges. The base and side covers are integrally molded as a one-piece structure made of an electrically insulating material. The base supports terminals, such as a pair of power terminals and a ground terminal, which have wiring ends extending back of the base. The terminals may have male contacts, such as a pair of power blades and a ground pin extending through suitable openings in the dead-front base, or they may have female contacts receiving male contacts passing through suitable openings in the dead-front base, or both types of contacts. There are locking portions molded integrally with the dead-front base and the side covers which engage when the connector is closed and the back ends of the side cover are secured to each other, which locking portions resist relative motion between the side covers and the dead-front base so as to supplement and even replace, if needed, the holding action provided by the web hinges. The electrical cord may be wired to the terminals in the course of folding the side covers to close the connector by having the wiring ends of the terminals in the form of insulation displacement slots opening toward one of the side covers. The intact insulated conductor wires of the cord are laid against the open ends of the terminal slots and the conductors are pushed into the slots, as the side cover facing the open ends of the slots is being closed, by a pusher shelf or shelves integrally formed with that side cover and suitably disposed with respect to the terminal slots. The terminal slots may be curved to be along an arc centered at one of the hinges so as to facilitate the insulation displacement connection. Their open ends may be staggered such that one of the cord conductors would start being pushed into its slot before another. The dead-front base may have barrier walls which surround the wiring parts of the terminals to keep them out of electrical contact with each other and, moreover, to prevent the terminal slots from being undesirably widened when the insulated conductors are being pushed into them. The locking means which resist relative motion between the side covers and the dead-front base when the connector is closed may include integrally formed means for locking, in a permanently closed position, the side cover which does not face the open ends of the terminal slots so as to use the permanently locked side cover as a shelf against which the cord conductors can be laid prior to closing the other cover so as to make the insulation displacement connection.\nThe objects of the invention include providing a high strength, one-piece housing for a dead-front electrical connector, providing a housing of this type which fully insulates the wiring area of the connector, providing a dead-front electrical connector which affords great ease of wiring and assembly, providing a dead-front connector which eliminates the variability of workmanship associated with screw-type terminal wiring and stripping of conductor ends for screw-type terminals, providing a connector having a housing made up of a dead-front base and web-hinged side covers which lock to the housing when closed so as to remain locked even in the absence of the holding action of the web hinges, and providing a dead-front connector which is particularly simple and inexpensive to make, assemble and use. These and other objects of the invention are discussed in more detail in the appended Detailed Description."}
{"text": "1. Field of the Invention\nThe present invention relates to a belt cleaner for removing the leavings of conveying material, e.g. coke, stuck on the surface of a conveyor belt during the return movement of the belt.\n2. Description of Related Art\nA prior art cleaner for a conveyor belt is disclosed in U.S. Pat. No. 3,841,470 wherein the belt cleaner includes a plurality of cleaner means disposed thereon in the transverse direction of the conveyor belt at the return running side thereof. Each cleaner means has a scraping portion thereof which intends to come into contact with the belt surface, while all of which are arranged in linear relationship. More specifically, said cleaner means which includes a scraper having said scraping portion formed with a tip member at uppermost end thereof is yieldingly supported by a resilient member so that the scraping portion thereof becomes in contact with the belt surface under pressure by the resilience of said resilient member.\nIt is known that an endless conveyor belt, which is fitted for running in a circle between a drive pulley and an idler pulley, is supported on the forward running side thereof by a plurality of guide rollers arranged at “trough” angles so as to be curved in an arcuate shape in cross section in order to carry the transferring material in a stable manner. In other words, on the forward running side of the belt where the material is transferred forward, the belt is supported so that the lengthwisely extending central region thereof can arcuately be deflected downward. The transferring materials are consecutively transferred generally on the central region of the belt, which will result in intensive wear on the central region of the belt surface rather than on both side regions of the same.\nBecause of the tendency to deflection (termed generally as a “trough” tendency), the belt is deflected arcuately upward on the return running side of the belt where the belt runs backward after discharging the transferring material therefrom. In such an arrangement of the known cleaner, the scraping portions of the cleaner means arranged in linear relationship are prevented from matching the belt surface when coming into close contact with the same. More specifically, there is a clearance between the scraping portions arranged linearly and the central region of the belt surface. The clearance becomes greater when the amount of wear on the central region of the belt surface is increased. Consequently, the leavings of material on the central region of the belt surface cannot positively be scraped off by the scraping portions of the cleaner. Since the leavings of material adhered to the belt are substantially greater in amount at the lengthwisely extending central region of the belt surface than the side regions of the same, some of the scraping portions disposed in the center get worn intensively in proportion to the amount of scraped material. This allows said clearance to become greater gradually.\nUnder the circumstances, an improved type of belt cleaner was developed by Nippon Tsusho Kabushiki Kaisha under the product name “U-Type Cleaner” to which some patents were granted as shown in the following patent documents.\nPatent document 1: JP-H2-50009B (corresponding to EP0289659B1)\nPatent document 2: JP-H5-65407B (corresponding to EP0338118B1)\nPatent document 3: JP-H7-20766B\nPatent document 4: JP-H11-35135A (corresponding to EP0891934B1)"}
{"text": "At present, prepumping to a rough vacuum (or backing vacuum as it is often called) is usually carried out with an oil-sealed rotary pump which is both lubricated and sealed with hydrocarbon or fluorocarbon oil. Some of the oil molecules are degraded and fragmented into smaller molecules during the operation of the rotary pump and these small hydrocarbon and fluorocarbon molecules exhibit a high vapour pressure relative to that of the oil before the latter was used in the pump. It is difficult to prevent these small molecules from passing back from the pump and entering the vacuum vessel where they contaminate all the surfaces of the vessel and its contents by coating them with an adherent oily film.\nIn the more modern electron microscopes, in which the specimen is not heated by the electron beam to the extent that it was in earlier types of electron microscopes, oil vapour arising from the oil lubricated pumping system condenses to form a contamination on the specimen, obscuring fine detail and reducing resolution by acquiring an electrical charge which deflects the electron beam. Furthermore, in the scanning transmission type of electron microscope which is now in use, the electron beam is produced from a tungsten tip by field emission, and the presence of any oil vapour in the surrounding vacuum will seriously affect the stability of the electron beam current.\nIn the technology of producing integrated silicon chip circuits in a high vacuum environment, the presence of any oil vapour is likely to render the chip inoperative because of the deposition of a thin oil film which may prevent good contact between layers and may insulate segments which are intended to be electrically connected.\nIn these fields, and others, oil lubricated pumps have continued to be used but elaborate systems have been developed for condensing out oil vapour or otherwise preventing it from reaching critical sites. One such system utilises a trap filled with pellets of alumina or zeolite, or a trap maintained at liquid nitrogen temperature, in the pumping line connecting the backing pump with the high vacuum pump. However, these traps are never completely effective in condensing out the oil vapour, so some contamination of the vessel with oily vapour always occurs.\nAt present, the only oil-free pumps capable of prepumping a vessel down from atmospheric pressure to fractions of a mm Hg are sorption pumps but the use of these is time-consuming and expensive. Sorption pumps usually consist of a stainless steel canister filled with zeolite pellets which, when cooled to liquid nitrogen temperature, have the ability to absorb most atmospheric gases. The canister is first heated and pumped with a backing pump (which needs to be fitted with an oil trap) to remove air from the zeolite pellets. It is then removed from the backing pump, connected to the vessel to be evacuated and then cooled to liquid nitrogen temperature, whereupon it begins pumping and continues to do so until the zeolite becomes saturated with air. The pump must then be disconnected from the vacuum vessel and reprocessed by heating and pumping and again cooling with liquid nitrogen. Sorption pumps were invented to provide oil-free prepumping of systems which are to be evacuated to a very high vacuum by oil-free pumps such as sublimation pumps, ionization pumps or cryopumps. Despite the cost of the liquid nitrogen used for cooling them and the inconveniences involved in processing them, they are widely used for such purposes.\nThose oil-free mechanical vacuum pumps which are commercially available are quite incapable of producing high vacuum. Two existing commercial pumps of this type employ split polytetrafluoroethylene (PTFE) sealing rings backed by a split, spring-steel band. The claimed performance of these pumps against atmosphere is 23 mm Hg (absolute) in one case and 124 mm Hg (absolute) in the other, and a clear limiting factor on performance is the split in the steel band which would allow a degree of air leakage. As backing pumps, the utility of these pumps in the applications discussed above is limited to prepumping prior to the use of sorption pumps.\nA further mechanical oil-free pump developed by the present applicant is disclosed in Australian Pat. No. 481072. This pump was found capable of producing high vacuum conditions without the use of lubricating and sealing oil but the vacuum which could be achieved was limited by difficulties in sealing against gas leakage into the working space of the pump and by the need to have valves which had to be subjected to gas pressure to open. The vacuum which could be produced in the high vacuum stage of a multi-stage pump was then determined by the pressure required to open an exhaust valve in the high vacuum stage of the pump.\nImprovements which have proven successful in meeting these difficulties are disclosed in Australian Pat. No. 516210 and in co-pending Australian patent application No. 68083/81. According to Pat. No. 516210 gas passes from a cylindrical working space above the piston to an annular working space below the piston by way of a gas transfer passage opening at the end face of the cylinder above the piston. Patent application No. 68083/81 discloses alternative sealing ring assemblies which have proven especially effective in enhancing the sealing of the cylinder."}
{"text": "1. Field of the Invention\nThis invention relates to a method of making a magnetic recording medium by an ion plating technique, more particularly, to a method of making a magnetic recording medium having excellent magnetic properties, especially excellent coercive force.\n2. Description of the Prior Art\nIt is known that conventional .gamma.-Fe.sub.2 O.sub.3, Co-doped .gamma.-Fe.sub.2 O.sub.3, Fe.sub.3 O.sub.4 and CrO.sub.2 powdery ferromagnetic alloys can be dispersed into an organic binder to form a coating-layer type magnetic recording medium. However, ferromagnetic metal thin films produced by electroplating, electroless plating, sputtering, vacuum deposition or ion plating have attracted special interest as non-binder type magnetic recording media, i.e., media that do not use a binder.\nAs one condition required for magnetic recording media for high density recording, it has been proposed to increase the coercive force of the magnetic layer by making the magnetic layer very thin. Therefore, non-binder type magnetic recording media have been sought by the art since the thickness thereof can easily be made ten times thinner than that of binder type magnetic recording media and they have a high saturated magnetic flux density.\nIon plating is a technique which was proposed by D. M. Mattox (see Japanese Patent Publication 8328/69 and U.S. Pat. No. 3,329,601). In ion plating, both a substrate and an evaporation source are disposed in an inert gas atmosphere kept at 10.sup..sup.-3 - 10.sup..sup.-1 Torr, and a DC voltage of 0.1 - 7.0 KV is then applied therebetween, the substrate and the evaporation source being connected to the negative and positive terminals, respectively, to thereby generate a glow discharge. The inert gas ions in the glow discharge impinge upon the substrate surface to clean it. After the substrate surface is cleaned, the evaporation source is heated to evaporate the metal to be deposited on the support, and the evaporated atoms are ionized in the plasma generated by the glow discharge. The ions are accelated by a strong electric field in the dark space area of the glow discharge surrounding the substrate, and have high energy and are deposited on the substrate.\nIn ion plating, a considerable amount of evaporation atoms are ionized and reach the substrate surface at high kinetic energy (several tens to several KeV) and both adhere to the substrate surface and cause etching due to the bombardment of the vaporized ions during the production of the thin film, whereby strong ferromagnetic metal thin films of high adhesion are obtained.\nSuch thin films are suitable for magnetic recording media which are employed under severe conditions, i.e., at high relative motion against a magnetic head.\nIn addition, the fabrication of thin films by means of ion plating can be adapted to mass-production because the production rate of the thin films is similar to that of the vacuum deposition technique, and, further, adhesion to the substrate is high as compared to the sputtering technique by which a thin film having a quality similar to the one produced by ion plating is obtained. In addition, the production of thin films by ion plating has the advantage that waste water generated by electroplating or electroless plating techniques is not a serious problem.\nHowever, it is difficult to obtain magnetic thin films with magnetic properties as are required for magnetic recording media, particularly in coercive force (Hc) and rectangular ratio (Br/Bm), for example, a coercive force above 300 Oe and a rectangular ratio above 0.65, by means of ion plating. Even if a metal such as Fe, Co, Ni, etc., or an alloy thereof is ion plated on a substrate to form a magnetic thin film, the coercive force is 400 Oe at most and the rectangular ratio is on the order of 0.4. Therefore, the necessary magnetic properties for use as a magnetic recording media cannot be obtained."}
{"text": "An icon for a trash-can or recycle bin on a computer desktop represents a folder that temporarily stores any files that are intended for deletion. Since a user who has deleted a file can change his/her mind and decide to recover the deleted file, many operating systems provide a trash-can function. In particular, some operating systems supporting graphical user interfaces (GUI), such as Apple Mac OS or Microsoft Windows, provide a trash-can icon, so that the user can easily delete a file by dragging and dropping the file to the trash-can icon."}
{"text": "1. Field\nThis specification relates to a mobile terminal, and particularly, to a mobile terminal capable of controlling contents displayed on a display unit, and a control method thereof.\n2. Background\nTerminals may be divided into mobile/portable terminals and stationary terminals according to their mobility. Mobile terminals (mobile device, portable device, portable terminal) can be easily carried and have one or more of functions such as supporting voice and video telephony calls, inputting and/or outputting information, storing data and the like.\nAs it becomes multifunctional, the portable terminal can be allowed to capture still images or moving images, play music or video files, play games, receive broadcast and the like, so as to be implemented as a multimedia player. Various new attempts have been made for the multimedia devices by hardware or software in order to implement such complicated functions.\nWith the improvement, the terminal may output contents on a display unit. However, when an event is generated while the contents are output, it causes inconvenience in that at least part of a screen for displaying the contents is obscured to display information relating to the generated event.\nWhen a user desires to view detailed information related to the generated event with a content display screen being output on the display unit, it causes user's inconvenience due to complicated manipulations."}
{"text": "Typically, stents are expandable, tubular metallic devices that are positioned within a patient's vasculature or other body lumen and expanded in order to support a vessel or body lumen and allow the flow of blood or other body fluids. Often, the stents are formed from a deformable metal and delivered using a balloon-type catheter. By advancing the catheter through the body lumen, the stent may be delivered to a desired position. Inflating the balloon then deforms the stent into an expanded configuration, seating it within the artery or other body lumen. Other implementations make use of a self-expanding stent formed from a psuedoelastic material that is delivered n a constricted condition and when released spontaneously expands to an enlarged configuration. In other embodiments, a stent made of shape memory alloy (NiTi) is inserted into the body lumen in a martensitic phase and transforms to austenite which has an expanded memory when raised to a temperature above the transformation temperature, usually less than 45° C.\nStents are often used in conjunction with an intravascular treatment for obstructive coronary artery disease. For example, ablation, atherectomy, balloon large percentage of percutaneous transluminal coronary angioplasty (PTCA) patients and rates can be even higher with other procedures. The prior art has employed a number of mechanical and pharmacological strategies to reduce the restenosis rate, but none have been particularly effective. Accordingly, stents have been proposed to maintain the patency of a treated vessel and prevent restenosis. Using stents, restenosis rates have fallen to less than 20%.\nRestenosis is thought to be a natural healing reaction provoked by injury from the intravascular procedure. The healing process frequently causes thrombosis and may lead to intimal hyperplasia that occludes the vessel. Although helpful in reducing restenosis, stents do not represent a complete solution. The framework of the stent may still allow migration and proliferation of the smooth muscle cells, while the stent itself can be thrombogenic. To address these problems, stents have been covered with DACRON, PTFE and autologous vein and the surface has seeded with endothelial cells or otherwise treated. Each of these solutions suffer from certain drawbacks, such as not being biocompatible, lacking sufficient mechanical strength, having a surface that is difficult to prepare, lack of ready availability and being thrombogenic.\nThus, there remains a need for a stent capable of minimizing restenosis while having a consistency similar to the native artery, a non-thrombogenic surface and sufficient mechanical strength as well as being biocompatible and readily available. This invention satisfies these and other needs."}
{"text": "The invention relates to a composite driving belt provided with a carrier and a plurality of transverse elements assembled slidably thereon."}
{"text": "A number of coffee drinkers have a concern about their level of caffeine consumption and many of these same consumers find that the aroma and/or flavor of the fully decaffeinated coffees presently available is unsatisfactory. These individuals are, therefore, likely to consume caffeinated coffee beverages but to minimize their intake of same. It is seen that a partially decaffeinated coffee beverage (say 40 to 60% decaffeinated) which approaches the flavor and aroma of fully caffeinated coffee beverages would enable the aforementioned consumers to increase their level of coffee consumption without increasing their intake of caffeine.\nA process of producing a partially decaffeinated soluble coffee product which contains a flavor and aroma improvement over present decaffeinated soluble coffee products would be an advantageous advance in the art, especially if the process could be accomplished without the use of organic solvents and without the necessity of constructing and operating separate coffee decaffeination facilities.\nIt has been found that attempts to produce partially decaffeinated coffee by merely reducing the percent decaffeination effected in known decaffeination processes does not result in coffee products which are significantly improved in flavor and aroma over their fully (i.e. 97% or more) decaffeinated counterparts. Decaffeination processes such as water decaffeination of green beans, as described in U.S. Pat. No. 2,309,092 to Berry et al.; the solvent decaffeination of green beans, as described in U.S. Pat. No. 3,671,263 to Patel et al.; or the decaffeination of aqueous extracts of roasted coffee, such as described in U.S. Pat. No. 2,933,395 to Adler et al. are sufficiently detrimental to coffee flavor and aroma that use of these processes to effect any significant degree of decaffeination is readily detectable. In addition, all these processes include the use of organic solvents such as the chlorinated hydrocarbons, which solvents are subject to governmental regulation.\nIt would, of course, be conceivable to use one of the conventional decaffeination processes to produce a decaffeinated coffee fraction which is then combined with a fully caffeinated fraction, which fully caffeinated fraction is able to provide a measure of desirable aroma and flavor to the final coffee product and beverage. This procedure, however, does not avoid the use of chlorinated or other solvents. Furthermore, to practice this procedure the coffee processor would be forced to either construct and operate a decaffeination unit or purchase decaffeinated green beans from an outside supplier.\nIt would be desirable if a soluble partially decaffeinated coffee could be produced which did not involve the use of organic solvents and which could be obtained from soluble coffee processing equipment without the necessity to construct separate decaffeination facilities. Published West German Patent Application No. 2,150,729 describes a process for producing a decaffeinated soluble coffee product with conventional percolation equipment; however, the process described therein is economically unattractive and the product produced thereby would not be regarded as desirable by caffeinated soluble coffee consumers."}
{"text": "The invention relates to cleaning processes of the type incorporating soluble oxidants and, in preferred embodiments, relates to removal of photo-resist (PR) and associated by-products during the cleaning stages of integrated circuit manufacture.\nIt is well known that gas solubilities in liquid media such as water can decrease with rising temperature. In particular, this has been problematic for cleaning systems based on dissolved ozone. When ozone functions as an oxidant in water, e.g., to effect removal of material from a surface, it can be counter productive to increase the intrinsic reaction rate of the ozone by increasing temperature. That is, when the intrinsic rate increases, the overall speed at which the system operates to remove material does not increase and may decrease, due to reduced solubility of the oxidant at elevated temperature.\nIn the semiconductor industry many manufacturing steps involve photolithography. Photoresist is commonly patterned and etched on an exposed surface of a partially fabricated semiconductor wafer in order to transfer a feature from a photomask to the surface. The feature, defined in the photoresist, is then transferred into the wafer material through, for example, a selective etch process. By way of illustration, this technique is commonly used to define zones for ion implantation, shallow trench isolation, polysilicon interconnect and Damascene trenches for metallization schemes.\nAfter the PR is patterned and etched, and after the patterned photoresist is used to etch the feature, it is necessary to remove photoresist or byproducts of the reactions. Conventionally, this has been accomplished with a cleaning process based on reaction of ozone with the material to be removed. For example, in a batch process, a group of wafers may be dipped in a solution of water and ozone, possibly under agitation, to effect the removal.\nIt is generally recognized that the reaction rate for ozone removal of PR and associated byproducts from a wafer surface is limited by diffusion of O3 to near the water/PR interface. Recently it has been observed that the rate of removal can be increased by first placing the wafers in a chamber of elevated temperature. With the wafers positioned above a pool of heated water, water vapor condenses upon the wafer surfaces while O3 is introduced into the chamber. In contrast to performing the reaction while the wafers are immersed in a bath of ozonated water, the wafers above the pool have a limited thickness of water on the surfaces. Since the PR strip rate is limited by diffusion of O3 through the DI water, providing a relatively thinner layer of water for a given partial pressure of O3 gas increases the net amount of O3 diffusion to the water/PR interface; and this will result in an increased reaction rate relative to the reaction rate which would result when the ozone diffuses into a bath of water to reach the PR.\nNonetheless, with both the solubility and the diffusion of the reactant species so limited, the approach of placing wafers over a pool of water to diffuse the species through condensate, at best, provides a rate-limited reaction. If the diffusion rate and overall rate of material removal could be further increased, significant economies may be made available. Specifically, PR removal in a single wafer processing system would become more cost effective.\nAccording to one embodiment of the invention, the removal of material in reactions having limited solubility and diffusion of gaseous material can be increased by controlling the thickness of the layer into which the reactant species diffuses. In another embodiment, the thickness of the layer and the diffusion of the gaseous material can be individually controlled. Apparatus and method are provided.\nA flow apparatus is provided for removal of material from a work piece having at least one reaction region containing the removable material. The apparatus may include first and second assemblies positionable in spaced-apart relation to form a zone extending between the two assemblies for movement of gaseous material. The first assembly may include a fixture portion positioned to receive the work piece with the reaction region of the work piece disposed in the zone to allow movement of the gaseous material thereover. A flow assembly is configured to transfer into the zone a gas comprising a condensable material and a reacting species.\nAccording to another embodiment a system is provided for removal of material from a workpiece having at least one surface region containing the removable material. The system includes a chamber, a flow component and a support apparatus. The flow component is configured to pass gaseous material into the chamber at a selectable rate and allow exit of liquid and gaseous materials from the chamber. The support apparatus is configured to position the workpiece thereon with the surface region of the work piece oriented for contact with gaseous material passed into the chamber by the flow component. The system may further include a thermal control configured to provide a differential temperature between a portion of a work piece positioned on the support apparatus and gaseous material passed into the chamber by the flow component.\nA method for chemical processing is also provided wherein the thickness of a layer of condensate on a surface is actively controlled and a reactant species is diffused from a gaseous region overlying the surface into the layer. According to an exemplary embodiment for removal of material from a work piece having at least one surface region containing the removable material, a workpiece is placed in an atmosphere comprising a condensable gas and a reacting species. The partial pressure of the condensable gas is controlled to limit the formation of liquid condensation on the surface region.\nAccording to a method for removing unwanted material from the surface of a semiconductor wafer, the wafer is placed in a zone having a controllable atmosphere and gaseous materials including a gaseous reactant and a condensing vapor are passed into the zone for contact with a wafer surface having the unwanted material thereon. The partial pressure of the condensing vapor in the zone is controlled to condense the vapor on the wafer surface at a selectable rate and cover the unwanted material with a condensate layer of desired thickness. The gaseous reactant is allowed to dissolve in the condensate.\nIn another embodiment for removing unwanted material from the surface of a semiconductor wafer, the wafer is placed in a chamber, gaseous materials including an oxidant and a condensing vapor are passed into the chamber for contact with a surface of the wafer having the unwanted material thereon, and vapor is condensed on the wafer surface at a predeterminable rate to cover the unwanted material with a layer of fluid. The oxidant is allowed to diffuse into the fluid at a selectable rate.\nThe foregoing background and summary have outlined general features of the invention. Those skilled in the art may acquire a better understanding of the invention and the preferred embodiments with reference to the drawings and detailed description which follow."}
{"text": "Electric power converter systems are used to transform and/or condition electrical power in a variety of applications. For example, electrical power converter systems may transform AC power from a power grid to a form suitable for a standalone application (e.g., powering an electric motor, lights, electric heater, household or commercial equipment, telecommunications equipment, computing equipment, and so forth). Another example standalone application is an uninterruptible power supply (UPS). UPS systems have become useful as backup power supplies for hospitals, financial institutions, industrial sites, telecommunications, and the like during interruptions of a public three-phase power supply grid. Increasingly, domestic homeowners also rely on UPS systems to supplement and/or replace power from the public power supply grid during failures.\nUPS systems typically incorporate some type of electrical power converter or other power transformation device/subsystem. An electrical power converter system may comprise one or more subsystems such as a DC/AC inverter, DC/DC converter, and/or AC/DC rectifier. Typically, electrical power converter systems include additional circuitry and/or programs for controlling the various subsystems, and for performing switching, filtering, noise and transient suppression, and device protection.\nBy way of example and historical explanation, power converter systems initially were built for specific applications. One early type of power converter was specifically designed for inverting direct current, constant voltage sources (e.g., batteries) to alternating current outputs (e.g., for operation of AC motors). Converters of this type are termed DC/AC “inverters,” and they have taken the form of, for example, transformers interconnecting a DC power supply with a plurality of logic control switches to generate the necessary alternating current waveform. A rectifier is another type of power converter for converting alternating current to direct current (AC/DC). Rectifiers have proven themselves useful for adapting household 110-volt alternating current to 12-volt direct current for operation of battery-powered appliances. Devices of this type have been as simple as a step-down transformer connected to a diode bridge and a smoothing capacitor for full-wave rectification. Direct current to direct current (DC/DC) converters have been provided for conditioning direct current power from a variable power source (e.g., a wind-driven direct current motor, photovoltaic panel or the like) for charging a battery or array of batteries.\nUninterruptible power supplies have been developed that permit power to be converted from a direct current power supply to a three-phase AC load in the event of a failure of the AC grid, and for recharging the DC power supply from the AC grid through the same apparatus when the AC grid is not in a failure mode. In other implementations, a UPS system transforms AC power from the AC grid into DC power, and then transforms the DC power into AC power, which is then provided to the AC load, instead of having the AC load directly powered by the AC grid. In the event of a failure of the AC grid, such a UPS system can continue to power the AC load by obtaining DC power from an alternate source (such as a battery) and then transforming that DC power into AC power.\nIn many situations, the AC load needs to be provided with AC power having a frequency (e.g., cycles per second) that matches the frequency of the AC grid. One reason to match frequency is that many devices comprising the AC load or devices coupled thereto have internal clocks or other components whose timing/cycles are based on the frequency provided by the AC grid. If the frequencies do not match, then such clocks can lose significant accuracy in their timing, ranging from a loss (or gain) of several seconds to even hours over a period of time.\nIn implementations where an AC load is directly powered from an AC grid, such frequency matching issues are less of a concern. That is, while the AC grid frequency may vary at times, the AC grid frequency always adjusts itself to average 60 Hz, for example, and can therefore be relied upon to provide accurate cycles.\nHowever, in implementations where the AC load is powered from a UPS rather than directly from the AC grid, the AC power provided by the UPS to the AC load may not necessarily match the AC grid frequency. In these situations, a frequency difference of 1% or greater, for example, from the AC grid frequency can cause noticeable cycle errors or other timing irregularities at the AC load. While oscillator tolerance in an internal microcontroller/microprocessor of the power inverter (or other power converter device) of the UPS can reduce cycle error, such reduction or other reliance on an internal clock of the microcontroller/microprocessor of the power inverter is not sufficient to guarantee substantially error-free cycles.\nClearly, minimum cycle error is needed by certain loads. Various institutions that rely heavily on computer data processing (such as financial institutions and telecommunication service providers) have little tolerance for significant cycle errors."}
{"text": "1. Field of the Invention\nThis invention relates to a method of making and testing integrated circuits, and a device used to perform such testing.\n2. Description of the Prior Art\nIntegrated circuits (ICs) comprise active and passive elements such as transistors, diodes, resistors, and capacitors, that are interconnected in a predetermined pattern to perform desired functions. The interconnections are effectuated by means of metallization layers and vias. A \"via\" is a hole through an insulation layer in which conductor material is located to electrically interconnect one conductive layer to another or to an active or passive region in the underlying semiconductor substrate. Present day technology generally employs two metallization layers that are superimposed over the semiconductor wafer structure. Integrated circuits and assemblies have become more complex with time and in a logic circuit, the number of integrated circuit logic units (ICLUs) and interconnects on a given size die have been substantially increased reflecting improved semiconductor processing technology. An ICLU can be a device (such as a transistor), a gate (several transistors) or as many as 25 or more transistors and other devices. As is well known in the art, these conductive contact points have a typical center-to-center spacing of about 6 to 15 microns (.mu.m).\nStandard processing to make logic structures (e.g., gate arrays) includes first fabricating as many as half a million transistors comprising a quarter of a million gates per die. Each semiconductor wafer (typically silicon but sometimes of other material such as gallium arsenide) includes many die, for example, several hundred. In one type of gate array, for example, the transistors are arrayed in rows and columns on each die, and each transistor is provided with conductive contact points (typically metal but sometimes formed of other conductive material such as polycrystalline silicon), also arrayed in rows and columns.\nIn the prior art, the next step is to use fixed masks to fabricate the conductive layers (sometimes called \"metallization layers\"), to connect together the individual gate-array devices. Typically two or sometimes three metallization layers are used.\nAfter this, the completed die is tested. If any of the devices on the die are defective, that die will fail an exhaustive test and be scrapped. Therefore, the more transistors per die the lower the manufacturing yield. In some cases redundant sections of a circuit are provided that can be substituted for defective sections of a circuit by fuses after metallization. Typically such redundant sections can be 5% to 10% of the total circuit."}
{"text": "The related background art known to Applicants, but which does not teach, disclose or suggest the present invention, includes the many prior patents showing side-by-side housing chambers disposed laterally of the mower path, and also the many prior patents showing a discharge duct extending upwardly and rearwardly of a mower housing for discharging air and grass cuttings blown therethrough by a rotating cutting blade in the housing."}
{"text": "1. Field of the Invention\nThe present invention relates to an improvement in a vehicular passive seat belt arrangement for selectively automatically restraining a vehicle occupant and releasing the occupant from restraint, and more particularly to a control system for such a passive seat belt arrangement to prevent its malfunction so as to improve safety for the vehicle occupant.\n2. Description of the Prior Art\nRecently passive seat belt arrangement has been increasingly mounted on automotive vehicles. A typical passive seat belt arrangement is configurated as follows: An upper end section of a seat belt is movably engaged with a guide rail extending in the fore-and-aft direction of the vehicle, so that the seat belt upper end section is selectively driven forward and rearward, i.e., toward its restraining position at which the vehicle occupant is restrained by the seat belt and toward its releasing position at which the vehicle occupant is released from restraint. Such driving of the seat belt is carried out in response to ON-OFF actions of a door switch adapted to be turned ON and OFF respectively when the door is opened and closed. This driving control manner is adopted upon paying attention to the fact that the door is necessarily opened when the vehicle occupant gets off the vehicle while closed during vehicle cruising, i.e., the seat belt arrangement is controlled t release the vehicle occupant from restraint upon the door switch being turned ON when the door is opened, while to restrain the vehicle occupant upon the door switch being turned OFF when the door is closed.\nHowever, with such a conventional control manner of the passive seat belt arrangement, even in case the door is opened contrary to the vehicle occupant's will, for example, owing to trouble of a door locking device, the door switch is turned ON and accordingly the seat belt upper end section is unavoidably moved to the releasing position at which the seat belt arrangement releases the vehicle occupant from restraint. This is because the driving control of the passive seat belt arrangement is made only in accordance with ON-OFF actions of the door switch. Such a control manner provides a possibility of the vehicle occupant being released from restraint with the seat belt arrangement during vehicle cruising, and therefore very dangerous for the vehicle occupant."}
{"text": "Plastic materials used in the manufacture of powder coatings are classified broadly as either thermosetting or thermoplastic. In the application of thermoplastic powder coatings, heat is applied to the coating on the substrate to melt the particles of the powder coating and thereby permit the particles to flow together and form a smooth coating.\nThermosetting coatings, when compared to coatings derived from thermoplastic compositions, generally are tougher, more resistant to solvents and detergents, have better adhesion to metal substrates, and do not soften when exposed to elevated temperatures. However, the curing of thermosetting coatings has created problems in obtaining coatings which have, in addition to the above stated desirable characteristics, good smoothness and flexibility. Coatings prepared from thermosetting powder compositions, upon the application of heat, may cure or set prior to forming a smooth coating, resulting in a relatively rough finish referred to as an \"orange peel\" surface. Such a coating surface or finish lacks the gloss and luster of coatings typically obtained from thermoplastic compositions. The \"orange peel\" surface problem has caused thermosetting coatings to be applied from organic solvent systems, which are inherently undesirable because of the environmental and safety problems that may be occasioned by the evaporation of the solvent system. Solvent based coating compositions also suffer from the disadvantage of relatively poor percent utilization; in some modes of application, only 60 percent or less of the solvent-based coating composition being applied contacts the article or substrate being coated. Thus, a substantial portion of solvent-based coatings can be wasted since that portion which does not contact the article or substrate being coated obviously cannot be reclaimed.\nIn addition to exhibiting good gloss, impact strength, and resistance to solvents and chemicals, coatings derived from thermosetting coating compositions must possess good to excellent flexibility. For example, good flexibility is essential for powder coating compositions used to coat sheet steel that is destined to be formed or shaped into articles used in the manufacture of various household appliances and automobiles, in the course of which the sheet metal is flexed or bent at various angles.\nFormation of a powder coating composition typically entails the dry mixing of flakes or granules of resin with the cross-linking agent and other ingredients, extruding the mixture at temperatures in the range of about 80.degree. to 130.degree. C., cooling the extrudate, and then chipping and grinding the resulting solid into particles of suitable size. This pulverizing operation typically produces a powder in which the particles are characterized by irregular shape and a broad size distribution."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display device and, more particularly, to a structure of a liquid crystal display device and its fabrication method.\n2. Background of the Related Art\nGenerally, a liquid crystal display device includes a lower substrate on which thin film transistor and pixel electrode are arranged, an upper substrate on which a color filter for displaying colors and a common electrode are formed. A liquid crystal is filled between those two substrates, and a polarizing substrate is installed on both sides of those two glass substrates and for pre-polarizing the visible ray (natural ray).\nFIG. 1 is a circuit diagram of the thus-structured liquid crystal display (LCD) device. The device includes a thin film transistor (TFT), namely, a switching circuit, a liquid crystal capacitor and a storage capacitor formed by the liquid crystal between the upper/lower substrate electrodes, a gate signal line and a data signal line.\nIf a voltage signal is applied to the gate line, the TFT turns on while the data voltage signal containing the picture image data is applied to the data signal line. The data voltage of the data signal line passes through the TFT and charges up the storage capacitor and the liquid crystal capacitor, thereby operating the LCD device.\nFIG. 2 is a plan view of the lower substrate of the related art LCD device. FIGS. 3A to 3H are sectional views illustrating the fabrication steps of the LCD device taken along lines I--I of FIG. 2.\nAs illustrated in FIG. 3A, a polycrystalline silicon is deposited on a transparent substrate 1 such as glass or quartz substrate, and patterned to form a semiconductor layer 2 on a predetermined portion of the substrate 1. The semiconductor layer 2 is used for an active area of TFT and an electrode of the storage capacitor.\nAs illustrated in FIG. 3B, a photoresist layer 3 is coated on the overall surface of the substrate and patterned to expose the semiconductor layer 2, which will be used as the lower electrode of the storage capacitor. Impurity ions, P or B, are implanted into the semiconductor layer 2 using the photoresist layer 3 as a mask.\nAs illustrated in FIG. 3C, the photoresist layer 3 is eliminated and a gate insulating layer 4 is formed on the overall surface of the substrate 1, including the semiconductor layer 2. A layer made of a material such as B or P doped polycrystalline silicon or silicide is deposited on the overall surface of the substrate 1, including the gate insulating layer 4, and patterned to form a gate line 5a and a common electrode line 5b. The common electrode line 5b serves as the upper electrode of the storage capacitor.\nAs illustrated in FIG. 3D, impurity ions, P or B, are implanted into the semiconductor layer 2 by using the gate line 5a as a mask, and the semiconductor layer 2 is heat-treated to activate the implanted ions and to form source and drain regions.\nAs illustrated in FIG. 3E, the first interlayer insulating layer 6, such as silicon oxide, is deposited on the overall surface of the substrate 1. The gate insulating layer 4 and the first interlayer insulating layer 6 are selectively removed to open a source region in the semiconductor layer 2, thereby forming a first contact hole 7.\nAs illustrated in FIG. 3F, a metallic layer is deposited on the overall surface of the substrate 1 including a first interlayer insulating layer 6. The metallic layer is patterned to form a data line 8 and is connected to the source region of the semiconductor layer 2 through the first contact hole 7.\nAs illustrated in FIG. 3G, a second interlayer insulating layer 9 is deposited on the overall surface of the substrate 1 including the data line 8. The gate insulating layer 4 and the first and second interlayer insulating layers 6 and 9 are selectively eliminated to expose a drain of the semiconductor layer 2, thereby forming a second contact hole 10.\nAs illustrated in FIG. 3H, the transparent conductive layer made of a material such as indium tin oxide (ITO) is deposited on the second interlayer insulating layer 9, and patterned to form the pixel electrode 11 which is connected to the drain region of the semiconductor layer 2. A passivation layer 12 is formed on the overall surface of the substrate 1 including the pixel electrode 11 to complete the formation of the lower substrate of the LCD device.\nThe related art device and its fabrication method have various disadvantages. For example, the sequential deposition of the semiconductor layer, the gate insulating layer and the common electrode line to form the storage capacitor decreases the aperture ratio of the LCD device. The common electrode line is too opaque and the aperture ratio is decreased by the size of the storage capacitor, namely, by 20-30% of the pixel area.\nFurther, since the gate line is formed by using a polycrystalline silicon including impurity ions, its resistance is too large for application to a LCD device of high image quality. Moreover, the hydrogen ions cannot enter into a channel area through the gate electrode during a hydrogenation process in which hydrogen radicals are fed into the semiconductor layer to enhance the performance of the device by the polycrystalline silicon TFT. Instead, the hydrogen radicals are laterally diffused through the gate insulating layer when the gate line is formed of a silicide material, such as WSix or MoSix. Such a process is time consuming to perform the hydrogen process the performance of the device.\nThe above references and/or description are incorporated by reference herein where appropriate for appropriate teachings of additional or alternative details, features and/or technical background."}
{"text": "Several conditions require sampling of blood. The sampling of blood can be used to determine the blood level of drugs with narrow therapeutic ranges, determine compliance or detect drugs of abuse. Blood sampling may also be necessary to determine adequate control of a disease state where, for example, without limitation, the patient can not regulate the level of a compound. For example, diabetics do not have the ability to control glucose levels in the blood to maintain such levels in a normal range. Glucose blood levels are needed to regulate the administration of insulin or other hypoglycemic or hyperglycemic agent.\nDiabetics need to test blood levels several times a day. However, the experience is not pleasant. There is pain associated with sticking a lance into the flesh in order to draw blood. Prolonged and repeated sticking causes skin abrasion as well. There is a psychological aversion to sticking a lance into oneself. Patient compliance for testing blood glucose levels is a problem.\nIt would be desirable to be able to test blood levels of biologically important compounds, such as glucose, or drugs without discomfort or the aversion of sticking a lance."}
{"text": "1. Field\nThe present disclosure relates to the field of ill conditioning, and specifically, determining whether square linear systems of equations are ill conditioned.\n2. Description of Related Art\nAnalytics systems are often used in business applications. Business Analytics use statistical and quantitative methods on data to perform forward looking analysis. Applying quantitative optimization methods to business analytics problems can enable companies to more effectively handle the complexity and large amounts of data in their operations.\nIBM® ILOG® CPLEX®, also known as CPLEX, named for the simplex method in combination with the C programming language, is software from the IBM Corporation. CPLEX can solve mathematical optimization problems arising from business analytics models as well as other scientific computing applications. All of CPLEX's algorithms come from the field of mathematical programming, and they involve computing solutions to square linear systems of equations. CPLEX can solve optimization models in numerous industries, including production planning in manufacturing processes, crew scheduling in the airlines industry, vehicle routing and delivery in the transportation industry, and computational biology in medicine. However, within analytics systems, such as CPLEX, ill conditioning in the square linear systems of equations can occur.\nIll conditioning, which is a concept of scientific computing and numerical linear algebra, is a situation in which a small change in an input value to a mathematical model can result in a large change to the computed solution or output value. For example, square linear systems of equations can be determined to be ill conditioned if a condition number is large, and it can be determined to be well conditioned if the condition number is small. Further, as the value of a condition number increases, the change in the computed solution or output value has the potential to increase relative to the change in input data.\nComputers have a limited number of bits to represent numeric values. The typical floating point system used by computers cannot represent all rational numbers exactly, and any system with a finite number of bits cannot represent all real numbers exactly. Therefore, scientific computing applications often cannot exactly represent the mathematical systems they model. In other words, these floating point systems have finite precision.\nA computer's machine precision represents the rounding error that can arise from the finite precision in the representation of such numerical values. The condition number of a square matrix corresponding to a linear system of equations provides a multiplicative factor of the change in input value that can manifest itself in the computed solution or output value. Specifically, consider a matrix A with m rows and m columns, and m vectors x and b of variables and data, respectively. For a given change in the input in either A or b, the condition number of A provides a measure of the change in output in the computed solution x of the square linear system Ax=b. The condition number of the matrix A associated with this linear system can be quantitatively derived asK(A)=∥A∥*∥A−1∥\nIf a computing application that solves such square linear systems is run on a computer with finite precision, then the machine precision provides a measure of the minimum change in the input. Letting Δ represent the small change in input to the data of the linear system in A or b, then K(A)*Δ provides an upper bound on the change to the computed solution x. From this product of K(A) and Δ we see that larger condition numbers imply larger potential changes in the computed solution. Such large condition numbers can make seemingly irrelevant changes to the input yield much larger changes in the output.\nEven if a mathematical model formulator does not change the data, finite precision computers can introduce small changes into square linear systems of equations, which could result in a large change to the computed solutions or output values. For example, if a user moves their program from a machine having a Microsoft Windows® operating system to a machine having an IBM AIX® operating system, the move can change the machine precision enough to significantly influence results if square linear systems of equations are ill conditioned.\nHowever, there currently does not exist any method for quantitatively assessing ill conditioning in a square linear system of equations. The current state of the art does not determine what constitutes a large or a small condition number when assessing the conditioning of a square linear system of equations. Also, there is no known method for addressing ill conditioning in algorithms, such as linear, mixed integer, quadratic, and nonlinear programming algorithms, that solve a series of interrelated square linear systems of equations.\nIn the current art, many programs do not address the issue of ill conditioning of square linear systems of equations. At most, in the event square linear systems of equations are determined to be ill conditioned, the current art would merely display a single condition number. The user would then have to determine whether the value of the condition number is sufficiently large to make the system ill conditioned. Further, when addressing only a single condition number from a single linear system of equations, the situation in which a sequence of condition numbers are generated, is not addressed. In some instances, examining a single condition number would give a partial view of ill conditioning for a particular problem."}
{"text": "1. Field of the Invention\nThis invention relates to a brake booster device which is provided in a brake system of, for example, an automobile and which is operated to increase an operating force (or control force) applied from a brake pedal, and transmit the increased operating force to a master cylinder.\n2. Description of the Related Art\nFIG. 5 is a diagram schematically illustrating the construction of an example of a conventional speed density type control system for an automotive internal combustion gasoline engine. As shown in this figure, an injector 2 for injecting fuel, and an ignition coil 3 for generating sparks are attached to an engine body 1. Moreover, an intake (or inlet) manifold 4 is connected to the engine body 1. Further, part of the gas flowing from the engine body 1 toward the intake manifold 4 is fed back to the intake manifold 4 through a recirculation pipe 5. An exhaust gas recirculation (EGR) valve 6 is provided at a midpoint on the recirculation pipe 5.\nA throttle valve 8, which is opened and closed mainly by operating an accelerator pedal (not shown), is provided in an air suction pipe 7 for introducing air from an air cleaner into the intake manifold 4. Further, a bypass pipe 9 for connecting upstream and downstream portions of the air suction pipe 7, which are upstream and downstream to the throttle valve 8 respectively, with each other is attached to the air suction pipe 7. Moreover, a bypass air valve 10 is provided in the bypass pipe 9. The injector 2, the ignition coil 3, the throttle valve 8 and the bypass air valve 10 are controlled by a control section (namely, a computer) 11.\nA brake booster body 14 for increasing operating force applied from the brake pedal 12 and transmitting the increased operating force to the master cylinder 13 of a brake unit is connected with the intake manifold 4. A negative-pressure (or vacuum) chamber 14a which is communicated with the intake manifold 4 through a negative-pressure suction pipe 15, and an atmospheric pressure chamber 14b into which atmospheric air is introduced when pressing the brake pedal 12, are provided in the brake booster body 14.\nFurthermore, a booster pressure sensor 16 which is operative to detect a pressure in the negative-pressure chamber 14a and to output a signal representing the detected pressure to the control section 11, is connected with the brake booster body 14. A manifold pressure sensor 17 which is operative to detect a pressure in the manifold 4 and to output a signal representing the detected pressure to the control section 11, is provided in the manifold 4. The amount of fuel injected from the injector 2 is calculated by the control section 11 according to a signal output from the manifold pressure sensor 17. Further, an atmospheric pressure sensor 18 for detecting atmospheric pressure is provided in the control section 11.\nFIG. 6 is a partially sectional diagram showing the construction of the brake booster body 14 in FIG. 5. Similar construction of a brake booster is illustrated in, for example, the May 1996 issue of \"Automobile Engineering\" (Jidosha-Kogaku). As shown in FIG. 6, a power piston 22 capable of moving in lateral directions (as viewed in this figure) is enclosed in a booster casing (or case) 21. A diaphragm 23 is securely fixed onto the power piston 22. Circumferential edge portions of the diaphragm 23 are fixed to the booster casing 21. Thus, the inner space of the booster casing 21 is partitioned into the negative pressure chamber 14a and the atmospheric pressure chamber 14b.\nA poppet valve assembly 24 is attached to the central portion of the power piston 22. The poppet valve assembly 24 is operated by the brake pedal 12 through a valve operating rod 25. Namely, when the brake pedal 12 is not depressed, the negative pressure chamber 14a and the atmospheric pressure chamber 14b are communicated with each other and are shut off from the outside. Further, when depressing the brake pedal 12, the atmospheric pressure chamber 14b is shut off from the negative pressure chamber 14a, while at the same time air is introduced into the atmospheric pressure chamber 14b.\nThe movement of the power piston 22 is transmitted to the master cylinder 13 through a push rod 26. The power piston 22 is pushed by a diaphragm return spring 27, which is disposed in the negative pressure chamber 14a, toward the atmospheric pressure chamber 14b.\nNext, the operation will be described. When the brake pedal 12 is not depressed, the negative pressure chamber 14a and the atmospheric pressure chamber 14b are communicated with each other and are shut off from the outside by the poppet valve assembly 24. Therefore, a pressure Pa in the negative pressure chamber 14a is equal to a pressure Pb in the atmospheric pressure chamber 14b (namely, Pa=Pb=a negative pressure).\nIn contrast, when depressing the brake pedal 12, the atmospheric pressure chamber 14b is shut off from the negative pressure chamber 14a and air is introduced into the atmospheric pressure chamber 14b. Consequently, Pa&lt;Pb=atmospheric pressure. Thus, a differential pressure between the internal pressure of the negative pressure chamber 14a and that of the atmospheric pressure chamber 14b is generated. Furthermore, the power piston 22 is moved toward the negative pressure chamber 14a against the force of the diaphragm return spring 27 by the differential pressure. Thus, the pushing force of the power piston 22 is transmitted to the master cylinder 13 through the push rod 26.\nThe pushing force of the push rod 26 due to the differential pressure is obtained by multiplying the differential pressure (Pb-Pa) by the area of the diaphragm 23. Thus, when depressing the brake pedal 12, the force obtained by adding the pushing force caused by the differential pressure to the brake pedal depressing force is transmitted to the master cylinder 13.\nThus, in the case of this conventional brake booster device, the operating force to be exerted on the master cylinder 13 is increased by utilizing the differential force between the (internal) negative pressure of the negative pressure chamber 14a and the atmospheric pressure. Consequently, there is the necessity of securing the aforementioned differential pressure at all times. Therefore, the internal pressure of the negative pressure chamber 14a is detected by the booster pressure sensor 16. In addition, the atmospheric pressure is detected by the atmospheric pressure sensor 18, and the differential pressure is monitored by the control section 11. Further, in the case that the differential pressure is insufficient, the throttle valve 8 or the air bypass valve 10 is closed. Thus, the internal pressure of the negative pressure chamber 14a communicated with the intake manifold 4 is reduced to thereby regulate the differential pressure.\nIn the case of the conventional brake booster device constructed as described above, since the internal pressure of the negative pressure chamber 14a is detected by the booster pressure sensor 16 which is communicated with the negative pressure chamber 14a, there is a need to provide the atmospheric pressure sensor 18 in the control section 11 to detect the differential pressure between the internal pressure of the negative pressure chamber 14a and the atmospheric pressure. This results in an increased manufacturing cost of the device. Moreover, there are variations among individual booster pressure sensors 16 and atmospheric pressure sensors 18. Consequently, measurement errors are increased.\nThis invention is provided to solve the aforementioned problems of the conventional device."}
{"text": "This invention relates to a coupling device that is easy to assemble, and just as easy to disassemble. This internal connecter device is shown here, as an example in the form of a rail. However it is not exclusively concerned with the formation of a rail. The xe2x80x9cclock lockxe2x80x9d can also be turned into the formation of a xe2x80x9cTxe2x80x9d-JUNCTION. Although its basic principles of operation do not change. With this method of assembly it is possible to achieve very flat external surfaces, at the point of connection, between two separate rails.\nA system for joining two rails includes a rotatable disc in a slot in an end of a rail, where the disc has a pair of diametrically opposed notches, a pair of diametrically opposed first sectors having a first radius, and a pair of diametrically opposed second sectors having a second radius larger than the first radius. A movable stopper in the rail is urged into the slot at a position where the stopper engages one of the first sectors when a pin in the rail engages one of the notches, and one of the second sectors is between the one first sector and the notch engaging the pin."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for treating a site contaminated with a hazardous substance. More particularly, the invention relates to a method and apparatus for continuously, and automatically, remediating soil and groundwater in a site contaminated with a hazardous substance.\n2. Description of the Prior Art\nWorldwide, the storage and transportation of hazardous substances requires countless tanks and pipelines. Substances typically processed in these facilities include petroleum distillates, industrial solvents, and industrial wastes. Due to the dangers presented by storing and transporting concentrated solutions of these materials, engineers and government agencies conduct extensive research and development to insure the tanks and pipelines used will effectively and safely contain these materials. Weathering and unforeseen engineering limitations, however, often cause the facilities to fail, resulting in the release of chemicals into the environment.\nThe failure of storage tanks and pipelines impacts on the physical and economic vitality of the contaminated area. The release of concentrated chemical solutions typically causes extensive damage to the local ecosystem by contaminating or killing indigenous plant and animal life. More remote ecosystems may be affected if the contamination migrates from the site by entering the local groundwater flow. Further, local laws often require that contaminated areas be remediated, and certified as such, prior to sale or rental. These laws often result in land being left fallow for years or decades. In extreme cases, the inability to remediate contaminated sites renders industrial or residential areas uninhabitable, forcing owners and employees to evacuate the area.\nSeveral methods are currently employed to remediate sites contaminated with hazardous substances. Unfortunately, available methodologies are often limited by soil conditions which, among other things, effect the rate of migration of hazardous contaminants. Also impacting on the usage of available remediation techniques is the size and configuration of the contamination plume. Finally, economic factors, particularly, labor and equipment costs, can, on occasion, render remediation procedures impossible to effectively commence, continue, or conclude.\nAt present, a common practice for the remediation of contaminated sites consists of locating a pumping system at or near the center of the plume. Contaminated groundwater is then pumped to the surface at a rate sufficient to create a cone of depression. The cone of depression causes floating contaminant to flow into the pumping well where the separation of contaminant from the groundwater can be accomplished using, for example, slow moving vertical flat, or round, belts made of oleophilic and hydrophilic materials. This is typically a labor intensive procedure, since the operation of the relevant equipment must be monitored on a virtually continuous basis in order to maximize product recovery. Thus, although effective for some applications, the physical size of these systems, their labor intensity, and power requirements, can render them inappropriate.\nA need has arisen for an improved method and apparatus for the remediation of sites contaminated with hazardous substances."}
{"text": "There are numerous methods for manufacturing gearwheels. In the case of cutting soft pre-machining, a differentiation is made between hobbing, gear shaping, generating planing, and power skiving. Hobbing and power skiving are so-called continuous methods.\nThe gear shaping method can be described or illustrated by a cylinder wheel gearing, since the intersection angle (also called the intersection angle of axes) between the axis of rotation R1 of the shaping tool 1 and the axis of rotation R2 of the workpiece 2 is 0°, as schematically shown in FIG. 1. The two axes of rotation R1 and R2 extend in parallel when the intersection angle of axes is 0°. The workpiece 2 and the shaping tool 1 rotate continuously about the axes of rotation R2 or R1, respectively, thereof. These rotational movements are shown by the reference signs ω1 and ω2. The shaping tool 1 performs a stroke movement in addition to the rotational movement, which is identified in FIG. 1 by the double arrow shx, and removes chips from the workpiece 2 during this stroke movement.\nSome time ago, a method was revived, which is referred to as power skiving. The fundamentals are approximately 100 years old.\nIn the case of power skiving, as shown in FIG. 2A, an intersection angle of axes Σ between the axis of rotation R1 of the power skiving tool 10 (also referred to as the skiving wheel) and the workpiece axis of rotation R2 of the workpiece 20 is specified, which is not equal to zero. The resulting relative movement between the power skiving tool 10 and the workpiece 20 is a spiral movement, which can be decomposed into a rotary component (rotational component) and a thrust component (translational component). A cylindrical screw gearing can be considered to be a drive technology analog, wherein the rotary component corresponds to the rolling and the thrust component corresponds to the sliding of the flanks. The greater the intersection angle of axes Σ is with respect to absolute value, the more the translational movement component required for the machining of the workpiece 20 increases. Specifically, it causes a movement component of the flank cutters of the power skiving tool 10 in the direction of the tooth flanks of the workpiece 20. In the case of power skiving, the sliding component of the meshing relative movement of the engaged gear wheels of the crossed helical equivalent gearing is utilized to execute the cutting movement. In power skiving, only a slow axial feed sax (also called axial feed) in parallel to the axis of rotation R2 of the workpiece 20 is required and the so-called impact movement is omitted, which is typical for gear shaping. Therefore, no reverse stroke movement occurs in the case of power skiving.\nThe cutting speed in power skiving is directly influenced by the rotational speed of the power skiving tool 10 or of the workpiece 20 and on the intersection angle of axes Σ used of the axes of rotation R1 and R2. The intersection angle of axes Σ and therefore the sliding component are to be selected so that an optimum cutting speed is achieved for the machining of the material at given rotational speed.\nThe movement sequences and further details of a power skiving method can be inferred from the already-mentioned schematic illustration in FIG. 2A. FIG. 2A shows the power skiving of external gear teeth on a cylindrical workpiece 20. The workpiece 20 and the tool 10 (a cylindrical power skiving tool 10 here) rotate in opposite directions, as can be seen in FIG. 2A, for example, on the basis of the angular velocities ω1 and ω2. The cylindrical tool 10 is inclined away from the workpiece 20 to generate kinematic clearance angles.\nFurther relative movements are added thereto. The above-mentioned axial feed s is necessary to be able to machine the entire gear tooth width of the workpiece 20 using the tool 10. The axial feed causes a displacement of the tool 10 in relation to the workpiece 20 in the parallel direction to the workpiece axis of rotation R2. The direction of this movement of the tool 10 is identified in FIG. 2A with sax. If spiral gear teeth are desired on the workpiece 20 (i.e., angle of inclination β2≠0), a differential feed sD is superimposed on the axial feed sax, which, as shown in FIG. 2A, corresponds to an additional rotation of the workpiece 20 about the workpiece axis of rotation R2. The differential feed sD and the axial feed sax are adapted at the calculation point AP to one another by computers such that the resulting feed of the tool 10 in relation to the workpiece 20 occurs in the direction of the tooth gap to be generated. In addition, a radial feed srad can be used, for example, to influence the crowning of the gear teeth of the workpiece 20.\nIn power skiving, the vector of the cutting speed {right arrow over (v)}c essentially results as the difference of the two velocity vectors {right arrow over (v)}1 and {right arrow over (v)}2, which are inclined in relation to one another by the effective intersection angle of axes Σeff, of the axes of rotation R1, R2 of tool 10 and workpiece 20. {right arrow over (v)}1 is the velocity vector on the circumference of the tool 10 and {right arrow over (v)}2 is the velocity vector on the circumference of the workpiece 20. The absolute value of the cutting speed {right arrow over (v)}c of the power skiving process can be changed by the intersection angle of axes Σ and the rotational speed in the crossed helical equivalent gearing. Then, as already mentioned, relatively slow axial feed sax only has a small influence on the cutting speed {right arrow over (v)}c in power skiving, which can be neglected. The axial feed sax is therefore not considered in the vector diagram having the vectors {right arrow over (v)}1, {right arrow over (v)}2, and {right arrow over (v)}c in FIG. 2A.\nFIG. 2B shows the power skiving of external gear teeth of a workpiece 20 using a non-inclined conical power skiving tool 10. In FIG. 2B, the intersection angle of axes Σ, the vector of the cutting speed {right arrow over (v)}c, the velocity vectors v, on the circumference of the tool 10 and {right arrow over (v)}2 on the circumference of the workpiece 20, and the angle of inclination β1 of the tool 10 and the angle of inclination β2 of the workpiece 20 are again shown. The angle of inclination β2 is not equal to zero here. The tooth head of the tool 10 is identified in FIG. 2B with the reference sign 4. The tooth face is identified in FIG. 2B with the reference sign 5. The two axes of rotation R1 and R2 do not intersect, but rather are arranged skewed in relation to one another. In the case of a conical power skiving tool 10, the calculation point AP is typically selected on the shared perpendicular line of the two axes of rotation R1 and R2, since tilting of the power skiving tool 10 is not required to provide clearance angles. The pitch circles of the spiral rolling equivalent gearing touch in the calculation point AP.\nA tool 10, which comprises at least one geometrically defined flank cutting edge, is used in power skiving. The flank cutting edge/flank cutting edges are not shown in FIG. 2A and FIG. 2B. The shape and arrangement of the flank cutting edges are among the aspects which have to be taken into consideration in practice in a specific design.\nIn the example shown in FIG. 2A, the power skiving tool 10 has the form of a straight-toothed spur gear. The external contour of the main body in FIG. 2A is cylindrical. However, it can also be conical, as shown in FIG. 2B.\nIn power skiving, the (tooth) gaps are typically generated in multiple cuts with respect to the tooth depth. This is therefore also referred to as a multiple cut strategy. The computer design of the power skiving tool is essentially performed in this case in view of optimum chip conditions for the final cut, during which the final flank surface is generated.\nThe individual cuts in different tooth depths are typically executed by a corresponding infeed of the power skiving tool in the tooth depth direction. In power skiving using a non-inclined power skiving tool, the infeed corresponds solely to a change of the axial distance between the axes R1 and R2.\nSeveral studies of power skiving using multiple cut strategy have shown that different chip or cutting conditions result for the individual cuts (at the different infeed depths). However, this is the case in particular if the same power skiving tool is used for the individual cuts in power skiving using multiple cut strategy. The chip or cutting conditions are different, since in particular the respective active pitch circle on the workpiece shifts in relation to the design position for the final cut.\nThe chip or cutting conditions even change significantly in this case, as precise simulations have shown. If inclined power skiving tools are used, in the extreme case, skiving cannot be performed without collision using an optimum power skiving tool for the final cut at a machining depth of approximately 25%.\nThese differences of the chip forming conditions at the various cutting depths not only result in wear of different levels on the power skiving tool, but rather they can also have an influence on the quality of the generated tooth flanks.\nFurthermore, it is characteristic for power skiving that the part of the chip removed from the incoming tooth flank eF is thinner than the part of the chip removed from the outgoing tooth flank aF. This method property of power skiving can also be problematic in the case of a multiple cut strategy."}
{"text": "Ostomy openings cannot be controlled at will, and in many cases patients with such surgically made openings or incontinent urine discharge openings are provided with a bag or pouch to collect the discharge from the openings. Many such pouches are disposable, i.e. for use only once, and the invention is not concerned with disposable pouches. In many other cases, however, it is preferred to use a drainable pouch provided with a discharge opening through which the pouch is emptied at more or less regular intervals. Many such pouches have been described, i.e. in the patent literature. The invention may be used in connection with any such known type of reusable collecting pouches. They are always provided with an inlet opening for the bodily excretions and around the inlet opening or elsewhere there is some device adapted to affix the pouch to the body of the patient. Moreover, there is frequently a venting aperture provided with a deodorizing filter to as to prevent bulging of the pouch due to intestinal gases. As the invention is not concerned with any such detail, they will not be discussed further.\nThe discharge opening of drainable pouches of the kind in question may be a slit in one of the pouch walls or, more commonly, merely an interruption of the seam joining the pouch walls together. As the main portion of the pouch is normally fairly wide, e.g. some 10 to 20 cm, it is usual to place the discharge opening in a narrower bottom extension of the pouch, e.g. having a width of some 4 to 6 cm. Moreover, it is common practice to close the discharge opening between the emptyings of the pouch by a clip, part of the narrower bottom portion being folded one or a few bends back and around the clip, the ends of which are doubled back over the folds of the bottom portion of the pouch. By experience such closing of the pouches is quite reliable.\nThe clips in question are usually strips made of a flexible yet relatively rigid and substantially unelastic material and in principle are wellknown, e.g. as closure means for all kinds of bags, e.g. for groceries or kitchen waste. They are frequently made of some rather rigid cellulosic material such as carton, or of some plastic material, and they are often reinforced with one or more metal threads.\nIn some cases the buyers of such drainable pouches use loose, independent clips. As these are easily lost in the process of emptying and possibly rinsing the pouch, it has been customary to glue or otherwise permanently affix the clip to the pouch adjacent the discharge opening thereof.\nHowever, this has the drawback that the clip and especially the narrow spaces and edge portions around it unavoidably become dirty during the emptying the pouch; and while the latter can be easily rinsed and cleaned because of the smoothness of the plastic sheet from which it is made, so is not the case with the clip and its immediate surroundings. The contamination of the clip and its close surroundings is unhygienic and unpleasant for the user, especially for patients with ileostomies and colostomies, for whom the problem is particularly pertinent. Moreover, it may render is difficult to keep the underwear clean.\nFrom FIG. 3 of the German published patent application DE 33 25 299 Al (having priorities from U.S. applications Nos. 398,913 and 491,256 of July 10, 1982, and May 3, 1983, respectively) there is known an ostomy pouch of the general type here concerned, having a narrowed bottom portion with the discharge opening right at the bottom of the bottom portion. A short distance above the discharge opening the bottom portion is provided with two lateral extensions, each forming a small pocket. The ends of a closing clip fits into and are placed in these pockets whereas the entire intermediate portion of the clip is free and situated at one of the walls of the pouch. Thus, when emptying the pouch, part of the waste material to remove may easily come into the space between the clip and the pouch wall. This can be avoided by removing the clip from the two pockets but firstly that in itself is unpleasant and may involve a risk of losing the clip; and moreover, waste material may easily come into the pockets in question, and it will be very difficult to rinse said small pockets since they have a width of the order of, e.g., 6-12 mm."}
{"text": "The present invention relates to a mechanical aparatus that extracts underground materials such as oil and so forth. It is a walking beam balanced span and moment regulating economized oil pump and a walking beam oil pump stroke regulating device, the latter of which is especially applicable to fore and rear walking beam oil pumps.\nThere are a variety of oil pumps which normally fall into the following categories: walking beam oil pumps and non-walking beam oil pumps in terms of the presence of the walking beam; conventional walking beam oil pumps, fore walking beam oil pumps, offset walking beam oil pumps and special-shaped walking beam oil pumps in terms of their structure; crank balance walking beam oil pumps, sidespin walking beam oil pumps, rotary walking beam oil pumps, suspended walking beam oil pumps and double-head walking beam oil pumps in terms of their mode of balance; and motor-driven and diesel-driven walking beam oil pumps in terms of the driving mode.\nAt present, a walking beam oil pump typically consists of an arc, a walking beam, a support, a pitman, a base, a crank, a decelerator, a power machine, and a string device. The support crank, decelerator, power machine and so on are fixed onto the base of the oil pump and the walking beam is hinged to the support via the central bearing. The pitman is hinged to the crank at one end, and to the walking beam through a bearing at the other end. The arc, on which is mounted the string device, is installed at the front end of the walking beam.\nWhen the oil pump is in operation, the load applied to the string device is called the arc suspension point load. The arc suspension point load varies around a maximum value during the upward oscillation or the up stroke of the arc, and the arc suspension point load varies in the vicinity of a minimum value during the downward oscillation or the down stroke of the arc. Therefore, the power machine is required to produce a large amount of energy during the up stroke whereas in the down stroke, the power machine has to do negative work due to the gravitation instead of producing a large amount of energy. Thus, in order to reduce such irregularity of load, balance devices are attached to oil pumps. However, although some of the oil pumps succeed in saving energy by means of the above-mentioned balance devices, there are many drawbacks with them, such as, complicated structures, poor balance performance, huge dynamic load, high energy consumption, low reliability, inconvenient operation and high costs of maintenance.\nIn particular, when some under-well technological parameters, such as the diameter of the oil pump, the depth of the deployment of the pump have been changed or when the viscosity of the oil and so on have varied, the difference between the maximum and minimum values of the oil arc suspension point load also varies accordingly. The output of the power machine has to be increased to adapt the oil pumps to such changes, which, however, inevitably increases the consumption of energy.\nCranes are usually needed to adjust the stroke of the conventional waking beam type oil pumps. This is achieved by hoisting the arc first with the crane, then the crank pin assembly connected to the pitman is taken out of the stroke regulating hole of the crank and then mounted in the stroke regulating holes of the crank to which the adjustment is to achieve. The stroke adjustment is not completed until the arc is deposited down and the crane has been driven away. Thus, a complicated series of procedures must be followed whenever adjustment of the stroke is needed, which costs a lot of time and labour, greatly influencing the efficiency of productivity.\nWith a view to the above problems, one object of the present invention is to provide a walking beam balanced span and moment regulating economized oil pump whose balance can be modified as the difference between the maximum and the minimum values of the arc suspension point load changes. The present invention may further adjust its balance along with the variation of the difference between the maximum and the minimum values of the arc suspension point load.\nAnother object of the present invention is to overcome the deficiencies of the conventional technology and to provide stroke regulating device applying to a walking beam oil pump which is simple in structure, easy to operate, time-saving and labour-saving.\nAccording to the first aspect of the present invention, the walking beam balanced span and moment regulating economized oil pump consisting of an arc, a walking beam, a support, a pitman, a base, a crank, a decelerator, a power machine, and a string device, is characterized in that the upper or lower part of the front end of the boom is hinged to the rear end of the walking beam while between the lower or upper part of the front end of the boom and the rear end of the walking beam is mounted a weighted arm regulator, and at the rear end of the boom is mounted a weighted device. The boom is made up of at least two segments, each of which is hinged upon each other at the upper part or the lower part of its front end, and between the lower part and the upper pat of the front end of each of the segments is mounted a weighted arm regulator. The weighted arm regulator consists of an adjustable pin roll and a mounting hole, and at the lower part of the front end of the boom there is a mounting hole whose diameter is larger than that of the adjustable pin roll, which is placed between the front end of the boom and the rear end of the waling beam through the hole. A seizure is fixed at the top of the boom corresponding to the mounting hole or at the end of the walking beam, and the seizure comprises at least one seizure notch. The weighted arm regulator consists of a wedge and a lead screw, in which a base plate of the is lead screw is fixed at the front end of the boom or the rear end of the walking beam, and the lead screw is seated on the base plate with the wedge fixed on it. The weighted arm regulator consists of a screw ejector pin and a screw ejector pin base, the latter of which is fixed at the front end of the boom or the rear end of the walking beam and the former mounted on the screw ejector pin base. The weighted means, which is hinged or articulated on the rear end of the boom, consists of an upper weighted box hinged or articulated on the rear end of the boom and at least one lower weighted box which can be mounted on the upper weighted box or removed from the upper weighted box. An unloading bumper bracket is mounted on the base corresponding to the lowest point of the weighted means. The boom is thick at the upper part and thinner at the lower part, thus constituting a uniform strength beam. The boom and the walking beam constitute either a trapezium or an angle or at, arc.\nAccording to the second aspect of the present invention, a walking beam stroke regulating device consisting of an arc, a wailing beam, a support, a pitman, a base, a crank, a decelerator, a power machine, and a string device, is characterized in that the fixing end of the stroke regulating strutting is positioned on the support or the walking beam, and the walking beam or the support corresponding to the strutting end of the stroke regulating strutting has a hinge housing, on which installed by the strutting end of the stroke regulating strutting when regulating the stroke. A guide sleeve is located at the fixing end of the stroke regulating strutting, and the fixing end is constituted by a flexible joint mounted together with the guide sleeve via restrictingnuts. The flexible joint is hinged to the hinge housing of the support or the base of the walking beam, that is, it is connected to the hinge housing of the support or the base of the walking beam with a rotable pin. A nut is placed at the strutting end of the stroke regulating strutting, and a turnbuckle is mounted on the strutting end through the screw nut. There is a restrictingscrew nut inside of the turnbuckle. The strutting end of the stroke regulating strutting is hinged on the hinge housing. The strutting end of the stroke regulating strutting is hinged on the hinge housing via a rotable pin. A split pin is mounted on the rotable pin. There is a hinged lever hole opened on the stroke regulating strutting."}
{"text": "The compound (E)-.alpha.-[2-n-butyl-1-[(4-carboxyphenyl)methyl]-1H-imidazol-5-yl]methyl ene-2-thiophenepropionic acid monomethanesulfonate is known by the name \"eprosartan\" and is the subject of U.S. Pat. No. 5,185,351 (the '351 patent), issued Feb. 9, 1993. This patent discloses in Example 41 a process for making the anhydrous form of (E)-.alpha.-[2-n-butyl-1-[(4-carboxyphenyl)methyl]-1H-imidazol-5-yl]methyl ene-2-thiophenepropionic acid monomethanesulfonate. Additionally, the '351 patent discloses conventional techniques for formulating (E)-.alpha.-[2-n-butyl-1-[(4-carboxyphenyl)methyl]-1H-imidazol-5-yl]methyl ene-2-thiophenepropionic acid monomethanesulfonate and Example; 108-111 specifically detail the preparation of certain formulations. This compound is claimed to have utility in blocking angiotensin II receptors and to be useful in the treatment of hypertension, congestive heart failure and renal failure.\nIt has been found that the monohydrate of (E)-.alpha.-[2-n-butyl-1-[(4-carboxyphenyl)methyl]-1H-imidazol-5-yl]methyl ene-2-thiophenepropionic acid monomethanesulfonate is formed during the vacuum drying of the dihydrated form of this compound or when the anhydrate of (E)-.alpha.-[2-n-butyl-1-[(4-carboxy-phenyl)methyl]-1H-imidazol-5-yl]methy lene-2-thiophenepropionic acid monomethanesulfonate is granulated with water, stored at 50.degree. C. overnight and vacuum dried overnight at ambient temperature."}
{"text": "Oscillator circuits are used in computers, computer peripherals, counters, timers, calculators, phase-locked loops, digital multimeters, oscilloscopes, and numerous other applications. An oscillator circuit may be used to provide a clock signal, to produce an accurate waveform, and the like.\nA type of oscillator circuit referred to as a relaxation oscillator generally operates as follows. The relaxation oscillator charges a capacitor with a reference current, and discharges the capacitor rapidly if the capacitor voltage reaches a reference voltage. The charging and discharging of the capacitor generates a sawtooth waveform that is used to generate the output signal. Alternatively, the polarity of the current may be reversed when the reference voltage is reached, generating a triangle waveform rather than a sawtooth waveform.\nAlso, the reference voltage is typically provided by a bandgap reference, so that the reference voltage is substantially independent of temperature and supply voltage. The current used to charge the capacitor is also typically substantially independent of temperature and supply voltage. In this way, the oscillator frequency can be substantially independent of variations in temperature and supply voltage."}
{"text": "In general, a window (opening) of a house or building is a location for the transfer in and out of a lot of heat. For example, a rate of heat loss flowing from a window in a case of heating in winter is approximately 48% and a rate of heat inflowing from a window in a case of cooling in summer is even approximately 71%. Therefore, it is possible to obtain an effect of enormous energy saving by appropriately controlling light or heat through a window.\nA light control glass has been developed for such a purpose and has a function of controlling an inflow and/or outflow of light and/or heat.\nThere are some kinds of methods for conducting a light control of such a light control glass. Among those, 1) a material with a light transmittance that is reversibly changed by applying an electric current or an electric voltage thereto is referred to as an electrochromic material, 2) a material with a light transmittance that is changed depending on a temperature is referred to as a thermochromic material, and further, 3) a material with a light transmittance that is changed by a control of an atmospheric gas is referred to as a gasochromic material. Among these, a study of an electrochromic light control glass that uses a tungsten oxide thin film for a light control layer is most advanced, so that a stage of practical application has generally been attained at present and a marketed product has also been provided.\nHowever, an electrochromic light control glass wherein a tungsten thin film is used for this light control layer is such that a principle thereof is that light is absorbed by the light control layer to conduct a light control. Therefore, there is a problem in that the energy saving effect is degraded because a light control layer absorbs light so as to be heated and it is also re-radiated into the room interior. In order to eliminate this, a light control is not conducted by absorbing light and it is necessary to conduct a light control by reflecting light. That is, a material (reflection-type light control element) has been desired that has a characteristic in such a manner that a state thereof is reversibly changed between a transparent state and a reflection state.\nFor a material that has such a characteristic, it has been found and reported in recent years that a state is reversibly changed between a transparent state and a reflection state due to a hydrogenation and a dehydrogenation of a rare earth metal such as a yttrium or a lanthanum (see, for example, U.S. Pat. No. 5,635,729).\nOtherwise, an alloy of a rare earth metal such as a gadolinium and a magnesium (see, for example, U.S. Pat. No. 5,905,590), an alloy of a magnesium and a transition metal (for example, see U.S. Pat. No. 6,647,166), and an alloy of an alkaline-earth metal such as a calcium and a magnesium (for example, see Japanese Patent Application Publication No. 2010-066747) have already been known as a material that has a reflection-type light control characteristic (light control mirror characteristic).\nHowever, a switching repetition durability of the above-mentioned reflection-type light control element between a transparent state and a reflection state is low. Accordingly, a method that interposes a buffer layer between a layer that has a reflection-type light control characteristic and a catalyst layer and further forms a hydrogen-permeable and water-repellent protection layer on a surface of the catalyst layer has been conducted in order to improve durability. However, there is a problem in that a switching repetition durability of about 1600 times is not provided even if a buffer layer or a protection layer is provided (for example, see Bao, et al., Solar Energy Materials & Solar Cells, Vol. 93, 1642 (2009))."}
{"text": "This invention relates to rail grinders and more particularly to rail grinders mounted on a vehicle that travels along the rails as they are being ground.\nAs railroad rails wear, crests and troughs are formed which define irregular waves extending longitudinally off the tracks. When the length of such waves is greater than the diameter of the circular grinding wheels used to smooth the rails, the rail grinding apparatus will follow the waves; this results in as much material being ground from the troughs as is ground from the crests. Though attempts have been made to solve this problem, the prior grinding apparatus followed the waves because the grinding discs were not mounted for movement that is essentially completely independent of the vertical movement of the rail carriage."}
{"text": "1. Field of the Invention\nThe present invention relates to surgical devices and methods to guide instruments that prepare the surface of bones and other tissues for implants that replace a damaged, diseased, or otherwise painful spinal facet joint.\n2. Description of Related Art\nTraumatic, inflammatory, metabolic, and degenerative disorders of the spine can produce debilitating pain that can have severe socioeconomic and psychological effects. One of the most common surgical interventions today is arthrodesis, or spine fusion, of one or more motion segments, with approximately 300,000 procedures performed annually in the United States. Clinical success varies considerably, depending upon technique and indications, and consideration must be given to the concomitant risks and complications. For example, Tsantrizos and Nibu have shown that spine fusion decreases function by limiting the range of motion for patients in flexion, extension, rotation, and lateral bending. Furthermore, Khoo and Nagata have shown that spine fusion creates increased stresses and, therefore, accelerated degeneration of adjacent non-fused motion segments. Additionally, pseudoarthrosis, as a result of an incomplete or ineffective fusion, may reduce or even eliminate the desired pain relief for the patient. Finally, the fusion device, whether artificial or biological, may migrate out of the fusion site.\nRecently, several attempts have been made to recreate the natural biomechanics of the spine by use of an artificial disc. Artificial discs provide for articulation between vertebral bodies to recreate the full range of motion allowed by the elastic properties of the natural intervertebral disc that directly connects two opposed vertebral bodies.\nHowever, the artificial discs proposed to date do not fully address the mechanics of motion of the spinal column. In addition to the intervertebral disc, posterior elements called the facet joints help to support axial, torsional and shear loads that act on the spinal column. Furthermore, the facet joints are diarthroidal joints that provide both sliding articulation and load transmission features. The effects of their absence as a result of facetectomy was observed by Goh to produce significant decreases in the stiffness of the spinal column in all planes of motion: flexion and extension, lateral bending, and rotation. Furthermore, contraindications for artificial discs include arthritic facet joints, absent facet joints, severe facet joint tropism or otherwise deformed facet joints, as noted by Lemaire.\nU.S. Pat. No. Re. 36,758 to Fitz discloses an artificial facet joint where the inferior facet, the mating superior facet, or both, are resurfaced.\nU.S. Pat. No. 6,132,464 to Martin discloses a spinal facet joint prosthesis that is supported on the posterior arch of the vertebra. Extending from this support structure are inferior and/or superior blades that replace the cartilage at the facet joint. The Martin prosthesis generally preserves existing bony structures and therefore does not address pathologies that affect the bone of the facets in addition to affecting the associated cartilage. Furthermore, the Martin invention requires a mating condition between the prosthesis and the posterior arch (also known as the lamina) that is a thin base of curved bone that carries all four facets and the spinous process. Since the posterior arch is a very complex and highly variable anatomic surface, it would be very difficult to design a prosthesis that provides reproducible positioning to correctly locate the cartilage-replacing blades for the facet joints.\nAnother approach to surgical intervention for spinal facets is provided in WO9848717A1 to Villaret. While Villaret teaches the replacement of spine facets, the replacement is interlocked in a manner to immobilize the joint.\nIt would therefore be an improvement in the art to provide a vertebral facet replacement device and method that provides a relatively high degree of mobility in the joint, while effectively removing the source of arthritic, traumatic, or other disease mediated pain with a minimum of patient discomfort."}
{"text": "Digital video streams typically represent video using a sequence of frames (i.e. still images). An increasing number of applications today make use of digital video stream encoding for purposes other than traditional moving pictures (such as movies and video clips)."}
{"text": "1. Field of the Invention\nThe present invention generally relates to image processing, and more particularly to a black level calibration method and system.\n2. Description of the Related Art\nUnless otherwise indicated herein, the approaches described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nImage sensors such as CMOS or CCD sensors are made up of an array of individual pixels, each of which collects photons incident on the image sensor. The number of photons collected in each pixel is converted into an electrical charge by a photodiode and this charge is then converted into an analog voltage, which may be amplified, adjusted, and converted to a digital value by an analog-to-digital converter, so that the information obtained from the individual pixels can be processed, usually by a digital signal processor, into a final digital image.\nMost image sensors require some form of calibration before use so that the data obtained from the image sensor can be used to produce digital images that faithfully reproduce the optical characteristics (e.g., intensity and color) of the scene or object whose image was captured. One type of calibration is referred to as black level calibration, which effectively sets a threshold below which digital data values obtained from the image sensor will be considered to represent a black level, or to represent the absence or substantial absence of light. Accurate black-level calibration helps to achieve a digital picture with full contrast and subtle details in dark shadow regions. If the black level is too low, information in dark areas may be lost. Conversely, if the black level is too high, signal range may be sacrificed.\nIn conventional systems, a border of an image-sensing array is surrounded with a number of rows and columns of light shielded, or black, pixels. These pixels provide black reference information or black pixel data to stabilize downstream image processing and establish the correct value for black in the output image.\nCalibration purely in the digital domain reduces the range of the system and reduces image quality. On the other hand, to accomplish high resolution and a wide calibration range simultaneously in the analog domain, existing solutions often involve circuits with large size and high power consumption."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit device and particularly to the technology that can be effectively applied to the measuring technique for semiconductor integrated circuit devices provided with a clock generating circuit for generating an internal clock signal corresponding to the clock signal which is supplied from an external terminal.\n2. Description of the Related Art\nThe search for references after application of the present invention has reported that an integrated circuit with PLL circuit provided with a jitter measuring circuit is disclosed in Japanese laid-open patent JP-A No. 2000-35463, an integrated circuit with PLL circuit provided with a detection circuit for evaluating jitter is disclosed in Japanese laid-open patent JP-A No. H6-104746, an integrated circuit with PLL circuit provided with a test circuit for calculating a lock time is disclosed in Japanese laid-open patent JP-A No. H11-2666, an integrated circuit provided with a frequency measuring circuit is disclosed in Japanese laid-open patent JP-A No. H11-23662 and an integrated circuit with PLL circuit provided with a frequency measuring circuit is disclosed in Japanese laid-open patent JP-A No. H9-197024.\nResearches are aided to improvement in the operation rate of a circuit with advancement of the semiconductor technology. For example, in a semiconductor integrated circuit device for data input/output operation synchronized with the clock signal of the frequency of 500 MHz or higher, an occupation rate of delay of signal in an input circuit for fetching, to the integrated circuit, the clock signal inputted from an external terminal for the clock period can not be neglected. An internal clock signal having compensated for signal delay in the input circuit can be formed with use of a PLL circuit or DLL circuit.\nCircuit elements formed in a semiconductor integrated circuit has comparatively large fluctuation in process. On the occasion of comprising a PLL or DLL circuit as explained above and forming an internal clock signal of high frequency as explained above, it can be thought that measurement to determine whether such clock signal should satisfy the desired frequency, phase or jitter characteristic or the like is now indispensable. Therefore, it can also be thought to use a high precision analog tester but in this case, an exclusive high-speed input/output circuit including the measuring terminals is required and moreover highly accurate determination becomes difficult even when an expensive measuring circuit is used because a signal delay in such input/output circuit becomes a serious problem. In addition, in the measuring circuit described in the above official gazettes, measurement accuracy is low and such measuring circuit is inadequate for PLL or DLL circuit which operates even in such higher frequency.\nAn object of the present invention is to provide a semiconductor integrated circuit device that has realized verification with higher accuracy of a plurality of operations for a clock generating circuit to form an internal clock signal. Another object of the present invention is to provide a semiconductor integrated circuit device that has enabled, through simplification of a structure, verification of various performances of an internal clock signal generating circuit. The abovementioned and other objects and novel characteristics of the present invention will be become more apparent from description of this specification and accompanying drawings thereof.\nTypical inventions among those disclosed in this specification will be outlined below. For the clock generating circuit to form an internal clock signal corresponding to an input clock signal inputted from an external terminal, a measuring circuit for measuring at least two or more items among the lock time until the desired internal clock signal corresponding to the input signal is obtained, the maximum frequency of the internal clock signal and the jitter of the internal clock signal is provided in order to execute, with higher accuracy, verification for operations of the clock generating circuit in the semiconductor integrated circuit."}
{"text": "Viral and bacterial infections are a major threat to public health. The emergence and expansion of life-threatening diseases caused by viruses and bacteria (e.g., acquired immune deficiency syndrome, severe acute respiratory syndrome, West Nile and Ebola hemorrhagic fevers, and tuberculosis), together with unmet conventional prevention approaches (e.g., vaccines) highlights the necessity of exploring new strategies that target these deadly pathogens.\nRecent studies have revealed that the TAM (Tyro3, Axl and Mer) tyrosine kinase receptors are positioned and function at a critical node of the innate immune response (Rothlin et al., Cell, 2007. 131(6):1124-36; Lemke and Rothlin, Nat. Rev. Immunol., 2008. 8(5):327-36). They are induced in dendritic cells (DCs) and macrophages by Type I interferon (IFN) receptors, which are themselves engaged as a consequence of toll-like receptor (TLR) activation upon encounter with pathogens (e.g., viruses and bacteria). The TAMs then act in concert with type IFN receptors to inactivate both the type I IFN receptors themselves, as well as the TLRs that initially trigger the inflammatory response to pathogens. In this way, the TAMs act as both pleiotropic inhibitors and integrated components of the innate immune response.\nThe innate immune system fights infection by viruses and bacteria in part through the production of Type I interferons (IFNs), a family of 13 alpha interferons and a single beta interferon (Borden et al., Nat. Rev. Drug Discov., 2007. 6(12):975-90). These agents display broad antiviral and antibacterial activities.\nGiven the foregoing, it would be desirable to have improved immunoenhancing agents, for instance for use in treating pathogen infections."}
{"text": "The present invention relates to a video camera, and more particularly to a portable video camera with a handgrip having a built-in action finder.\nThere are well known in the art various portable video cameras having a handgrip attached to the bottom of the camera body and a view finder mounted on the top of the camera body which is adjustable in angular position.\nThe recent tendency in portable video cameras is to provide a grip changeable in angle relative the camera body for easy handling. A problem in association with such a portable video camera is a difficulty of sighting a subject through a view finder of the portable video camera at a low angle camera position. This is due to the view finder being built in or mounted on the camera body. There are some large video cameras provided with adjustable view finders. Even with the provision of such an adjustable view finder, it is still difficult to handle the large video camera when using it at low angle camera positions because a grip of the large video camera is generally structurally integral with the camera body thereof.\nIn some instances, with video cameras of this type, it is troublesome to change the camera angle or position between a high angle and low angle camera position because it is necessary to bend the wrist or hand of an operator backward when shifting the camera in particular from a high angle camera position to a low angle camera position. Additionally, in some of the video cameras of this type two actions are necessary: adjusting the view finder, and thereafter the grip, in angular position. This brings the video camera to a standstill."}
{"text": "Powder metallurgy technology allows production of components which require high dimensional accuracy and have a complex structure in a shape markedly close to that of the finished product (in near net shape), thereby significantly decreasing the finishing cost. Therefore, many products produced by powder metallurgy are used as various components for machines and apparatuses in many fields.\nIn general, iron-based green compacts for powder metallurgy (green compacts) are produced as follows. First, an iron-based powder is mixed with alloying powder such as graphite powder and so forth, and lubricant powder such as stearic acid and lithium stearate to prepare an iron-based mixed powder. Then, the iron-based mixed powder is filled in a die, and is subjected to compacting, whereby the iron-based green compact is produced.\nThe iron-based powders are classified into iron powders (such as pure iron powder), alloy steel powder, and so forth, for example, based upon the components thereof. Also, the iron-based powders are classified into atomized iron powders, reduced iron powders, and so forth, for example, based upon the production method thereof.\nIn general, the iron-based green compacts are formed with a density of 6.6 to 7.1 Mg/M3. Furthermore, these iron-based green compacts are sintered to form sintered bodies. The sintered bodies are subjected to a sizing or a cutting process according to needs, whereby powder metallurgy products are produced. Furthermore, in some cases, the products are subjected to carburizing-quenching or bright-quenching after sintering for improving tensile strength or fatigue strength thereof.\nRecently, iron-based powder metallurgy products with high strength or high fatigue strength are strongly desired due to the development of components with reduced size and weight.\nIn general, alloying elements (Ni, Cu, Mo, W, V, Co, Nb, Ti, and so forth) are added to the iron-based powders for improving the strength of the powder metallurgy products.\nNote that examples of the methods for adding alloying elements include: a method for alloying the iron-based powder with a desired element (prealloying); a method for mixing an alloying powder (powder containing a desired alloying element) and the iron-based powder with or without binder; and a method for holding the mixture of the powder containing an alloying element and the iron-based powder at a high temperature so as to metallurgically combine these powders (diffusion bonding). Various properties of the alloy steel powder (or mixed powder), and various levels of uniformity and diffusion states of the alloying element after sintering are obtained depending on the method. Therefore, it is important to select the alloying element and the addition method for achieving the desired quality of the alloy steel powder (or mixed powder) or the desired quality of the sintered body.\nFor example, Japanese Examined Patent Application Publication No. 6-89365 discloses an alloy steel powder containing 1.5 to 20% by mass of Mo, which is a ferrite-stabilizing element, as a prealloy. According to the document in the sintering process of the aforementioned alloy steel powder, a single a phase is formed, leading to a high self-diffusion rate with respect to Fe. This accelerates sintering, resulting in a reduced size of the pores contained in the sintered body. Thus, pressure sintering of such an alloy steel powder provides a sintered body with improved densification. Furthermore, such an alloy steel powder contains no alloying element added by diffusion bonding, thereby providing a uniform and stable microstructure. However, the Mo content in the disclosure is relatively high, i.e., 1.8% by mass or more, leading to poor compressibility. This leads to the disadvantage that a green compact cannot be formed with high density (the density of the green compact). Accordingly, the sintered body obtained by performing a general sintering process (i.e., sintering in one step without pressurizing) has a low density, leading to insufficient strength and insufficient fatigue strength.\nOn the other hand, the pressure sintering method and the two-step sintering method including a repressing step have the disadvantage of high costs. Accordingly, a sintered body is preferably produced with high strength and high fatigue strength without involving such special sintering methods.\nOn the other hand, Japanese Examined Patent Application Publication No. 7-51721 discloses a steel powder which contains 0.2 to 1.5% by mass of Mo and 0.05 to 0.25% by mass of Mn as prealloyed elements, and which has a relatively high compressibility in compacting. However, it has been revealed by the present inventors that a single a phase is not formed using the aforementioned steel powder due to the Mo content of 1.5% by mass or less. Accordingly, the enhanced sintering between particles is not accelerated in a sintering step at a temperature (1120 to 1140° C.) of a mesh belt furnace generally used for powder metallurgy, leading to a problem of low strength of the sintering neck.\nWhile the Japanese Examined Patent Application Publication No. 7-51721 discloses an iron powder as a comparative example, which contains Ni (3.8% by mass), Mo (0.5% by mass), and Cu (1.4% by mass) by diffusion bonding, the Patent document describes that the iron powder has poorer strength than that of the aforementioned alloy steel powder disclosed as an invention in the Patent document.\nOn the other hand, Japanese Examined Patent Application Publication No. 63-66362 discloses a technique in which Mo is added to an iron powder as a prealloyed element so long as compressibility is not impaired (Mo: 0.1 to 1.0% by mass), and Cu and Ni are bonded on the surfaces of the iron particles in the form of a powder by diffusion bonding. This technique provides both preferable compressibility during the compacting and high strength after sintering. However, the aforementioned technique has a limited ability to improve tensile strength and fatigue strength by adding Cu and Ni since the iron powder containing Mo as a prealloyed element cannot be sintered sufficiently, as with the technique disclosed in the Japanese Examined Patent Application Publication No. 7-51721.\nOn the other hand, Japanese Unexamined Patent Application Publication No. 8-49047 discloses an alloy steel powder limiting Mn content to 0.3% or less by mass as a prealloyed element as well as containing Mo of 0.1 to 6.0% by mass and V of 0.05 to 2.0% by mass (as prealloyed elements). The aforementioned alloy steel powder provides a sintered body with high strength after heat treatment while maintaining the compressibility thereof. Also, the patent document discloses that the alloy steel powder may contain one or more kinds of elements of Mo (4% by mass or less); Cu (4% by mass or less); Ni (10% by mass or less); Co (4% by mass or less); and W (4% by mass or less) in the form of powders by mixture or diffusion bonding.\nOn the other hand, Japanese Unexamined Patent Application Publication No. 7-233401 discloses an atomized iron powder (alloy steel powder) which contains Mn of 0.03 to 0.5% by mass and Cr of 0.03 to less than 0.1% by mass as prealloyed elements. The aforementioned atomized iron powder having excellent machinability of the sintered body, as well as providing superior dimensional-accuracy thereof. Also, the aforementioned Patent document discloses examples of strengthening elements that can be used as prealloyed elements, which include: Ni (4.0% by mass or less); Mo (4.0% by mass or less); Nb (0.05% by mass or less); and V (0.5% by mass or less). Furthermore, the Patent document discloses examples of strengthening elements (alloy powders) that can be added by diffusion bonding, which include: a Ni powder (5.0% by mass or less); a Mo powder (3.0% by mass or less); and a Cu powder (5.0% by mass or less).\nHowever, according to the aforementioned techniques, such alloys are not designed from the perspective of the fatigue strength of components produced by sintering. This leads to difficulty in producing sintered metal components which satisfy the high fatigue strength desired in recent years, using a general sintering step.\nFor example, Japanese Unexamined Patent Application Publication No. 6-81001 and Japanese Unexamined Patent Application Publication No. 2003-147405 disclose alloy steel powders designed for improving the fatigue strength.\nThe Japanese Unexamined Patent Application Publication No. 2003-147405 discloses an alloy steel powder in which 0.5 to 1.5% by mass of Mo is bonded on the surfaces of a steel powder containing Ni of 0.5 to 2.5% by mass and Mo of 0.3 to 2.5% by mass as prealloyed elements by diffusion bonding. The aforementioned Patent document also discloses that a sintered body formed of the aforementioned alloy steel powder exhibits superior rolling contact fatigue strength after carburizing-quenching.\nOn the other hand, the Japanese Unexamined Patent Application Publication No. 6-81001 discloses an alloy steel powder in which Ni (0.5 to 5% by mass) and/or Cu (0.5 to 2.5% by mass) are bonded to an iron-based powder containing Mo of 0.05 to 2.5% by mass and at least one element of V, Ti, and Nb as prealloyed elements by diffusion bonding. The aforementioned Patent document discloses that the alloy steel powder provides a sintered body having superior rolling contact fatigue strength after carburizing-quenching, as well."}
{"text": "1. Field of the Invention\nThis invention relates to analog to digital converters and more particularly to converters which will operate at high speeds and over wide ranges of temperatures and other conditions to produce highly accurate and reliable results, without requiring self-calibration procedures. Converters of the invention have a relatively simple architecture and economically manufacturable.\n2. Description of the Prior Art\nHigh speed A/D converters have heretofore been provided which have included \"flash\" converter integrated circuits. Such circuits have included a plurality of voltage comparators which compare the level of an analog input voltage with a series of reference voltage levels which are supplied in steps of uniform magnitude, typically by a string of resistors which are connected to a reference voltage source. The magnitude of the input analog signal is indicated by the number of comparator output signals produced, i.e., by the number of reference voltage levels which are exceeded by the input voltage. A \"thermometer\" code output is thus produced which is converted to binary by a decode circuit.\nA problem with such flash converter circuits is that the number of required comparators and voltage levels doubles with each additional bit of the required resolution, so that an excessively large number of comparators and associated voltage levels would be required to obtain high resolution. In view of this problem, subranging A/D architectures have been used. A main range A/D converter is used which is of the flash type and which has a limited number of comparators, e.g., 32 or 64. The main A/D converter produces, in effect, a first approximation of the magnitude of an input analog signal. The first approximation so produced is then converted to an analog signal with a high speed precision D/A. The analog signal developed at the output of the D/A is then subtracted from the analog input signal, the result of the subtraction being applied to a second or \"subrange\" A/D converter to produce a binary coded output which is logically combined with a binary coded output from the main range A/D converter to produce the final binary coded digital output.\nIn an attempt to obtain high speed operation with prior subranging architectures, the analog input signal has been delayed before subtraction of the output signal of the high speed precision D/A therefrom, for the purpose of matching the delays of the main range A/D. A sample and a hold circuit is typically used as an input stage to apply a signal to the main range A/D and to obtain the required delay, a second sample and hold circuit may be used or, in the alternative, a delay line may be used. There is also a problem with digital skew going into the D/A and a register is oftentimes placed at the output of the main A/D to extend the time that the digital data is valid.\nPrior subranging architectures have had problems with respect to accuracy and reliability, especially with respect to the use of a delay line which adds an error source since termination of a delay line is never exact and multiple reflections may be produced to degrade the conversion accuracy when sampling high frequency signals. In addition, the input and output terminations on the delay line may attenuate the signal by a factor of two so that the gain of the sub-range amplifier must be increased by a factor of two. Offsets by buffers in the analog delay line path coupled with the extra sub-range gain required makes the architecture sensitive to temperature variations. In addition, operation at frequencies other than the frequency for which the converter circuitry was designed can cause large conversion errors.\nA specially constructed A/D/A (Analog/Digital/Analog) integrated circuit component has been tried in conventional architectures in an attempt to avoid problems, but with only limited success. The A/D/A Component includes comparators which drive decoding logic to develop a binary coded digital output and which also directly drive current sources to develop an analog output, thereby forming an A/D/A component in which the speed of development of the analog output is increased. It is found that although this component has important advantages, the use of this component with the conventional subrange architecture is still limited in providing a high speed and high resolution converter."}
{"text": "1. Field of the Invention\nThe present invention concerns a solenoid valve with programmable electronic control device, particularly for watering systems.\n2. Description of the Related Art\nThere are also known solenoid valves of the aforesaid type in which the control solenoid is made up of two parts, as described in U.S. Pat. No. 6,352,238, where a first part of the solenoid includes a piston for the operation of a hydraulic valve and a magnet for attraction of the same piston and is screwable on the external body of the solenoid valve and a second part includes an electric solenoid for the excitation of the magnet that is provided with xe2x80x9cfaston(copyright)xe2x80x9d type electric connectors and which is fastened to the first part in a removable way by means of a removable and re-insertable locking element.\nFor some uses it can be desirable to equip the single solenoid valve with its own electronic control device of programmable type, that makes the same solenoid valve independent from the general program of the watering system or that even prevents the need of having a control unit for the system.\nHowever a solenoid valve provided with an independent control device inevitably becomes object of thefts and vandalism acts, particularly during the periods of absence of the user.\nIn addition the electronic control device inset in the solenoid valve can deteriorate for infiltration of humidity or rain or other atmospheric events, in particular during the autumn and winter periods, when besides the watering is no longer required.\nFinally the programming operation of the electronic control device inset in the solenoid valve requires the intervention of the user or the programmer directly on the site where it is used, often inconvenient and difficultly accessible.\nAll that considered, object of the present invention is to realise a solenoid valve with autonomous programmable electronic control device that is free of the aforementioned disadvantages.\nAccording to the present invention, such object is reached by means of a solenoid valve comprising a hydraulic valve, a valve body in which said hydraulic valve is housed and a control solenoid for said valve, characterised in that it comprises in addition a programmable electronic control device that is housed together with said solenoid in a removable way inside a watertight openable container that is sealingly fastened to said valve body.\nOwing to the present invention a solenoid valve with its own electronic control device is thus realised that, when in use, preserves the same device as regards infiltration of humidity or rain or others atmospheric events and at the same time it allows, when desired, to extract the control device so as to have it protected from thefts and vandalism acts and to allow its convenient programming at home or in a laboratory.\nPreferably the solenoid is of the type that is made up of a first part including a piston for the operation of the hydraulic valve and an attraction magnet for said piston and of a second part fastened to the first one in a removable way and including an electric solenoid for the excitation of the magnet.\nThe electric solenoid too, that it is the most expensive part and subject to failure of the solenoid, can therefore take advantage of the favourable effects of the watertight containment and removability as the ones previously described for the electronic control device."}
{"text": "This specification relates generally to the field of loudspeakers. More particularly, this specification relates to a suspension system for the voice coil assembly of a loudspeaker.\nIn general, a loudspeaker includes a frame, a moving assembly, and a suspension system, which mechanically couples the moving assembly to the frame in a manner that permits the moving assembly to move relative to the frame. The moving assembly includes a diaphragm, which vibrates to radiate pressure waves that are perceived as sound. The suspension system, which may include one or more suspension elements, preferably permits movement along a single axis, so that contact between the moving assembly and the frame, or elements rigidly coupled to the frame, are avoided.\nOne common type of loudspeaker is a moving coil loudspeaker. In a moving coil loudspeaker, the moving assembly includes a voice coil assembly. The voice coil assembly includes a voice coil former, which is typically a tube with a circular cross section, but which may have some other form of cross section, such as square or rectangular. The moving assembly also includes a voice coil, which is typically formed by tightly winding an electrically conductive wire around the voice coil former. The diaphragm is mechanically coupled to the voice coil assembly. Audio signals, in the form of electrical current in the voice coil, interact with the magnetic field of a magnet assembly which is rigidly coupled to the frame, to cause the diaphragm to vibrate, radiating pressure waves that are perceived as sound.\nAnother type of loudspeaker is a moving magnet loudspeaker. In a moving magnet loudspeaker, the voice coil assembly is rigidly coupled to the frame, and the moving assembly includes a magnet assembly, mechanically coupled to the diaphragm. Audio signals, in the form of electrical current in the voice coil, interact with the magnetic field of the magnet assembly, to cause the diaphragm to vibrate, radiating pressure waves that are perceived as sound. In this specification, the examples are moving voice coil loudspeakers, but the principles described herein may be applied to moving magnet loudspeakers provided the suspension system has adequate properties such as lateral stiffness.\nA first common type of moving voice coil loudspeaker includes a voice coil assembly in which the diameters of the voice coil former and the diaphragm are substantially the same. In these loudspeakers, the outer-most edge of the diaphragm is attached to the upper periphery of the voice coil former. The moving assembly is typically secured to the frame of the loudspeaker by at least a first support element, commonly referred to as a “surround”, which has an inner edge secured to the moving assembly and an outer edge that is secured to the frame. Alternate embodiments may include a second support element, commonly called a “spider”, which includes an inner edge secured to a bottom portion of the voice coil former and an outer edge that is secured to the frame of the loudspeaker. This type of construction is typically found in smaller loudspeakers, such as tweeters, and possibly mid-range speakers.\nA typical issue encountered with smaller-sized loudspeakers is that as the loudspeaker becomes smaller, achieving low frequency response becomes more difficult. Low frequency response requires a loudspeaker to displace a larger volume of air to achieve the lower frequencies, as compared to achieving higher frequencies. The volume of air that a loudspeaker can displace is dependent upon the area of the diaphragm and the peak-to-peak excursion of the voice coil assembly that is allowed by the suspension. As the axial stiffness of the suspension is reduced to allow a greater excursion of the voice coil assembly, the radial stiffness of the suspension is usually similarly decreased. To operate at maximum efficiency, the suspension system in smaller loudspeakers should allow a maximum amplitude of axial displacement while constraining the voice coil assembly from moving side to side in order to avoid contact between the voice coil assembly and the other portions of the loudspeaker. As would be expected, as the stiffness of the suspension is reduced, greater side to side motion of the voice coil assembly is usually allowed. This is especially true in those loudspeaker embodiments that only include a single support element securing the voice coil assembly to the frame of the loudspeaker.\nIn a second type of moving coil loudspeaker, the moving assembly includes a diaphragm that is formed by a cone (or some other planar or non-planar surface, such as a concave surface) having a diameter that is greater than the diameter of the voice coil former. In this type of loudspeaker, an inner periphery of the cone is typically secured to an upper portion of the voice coil former and a first support member has an inner edge secured to an outer periphery of the diaphragm and an outer edge that is secured to the frame of the loudspeaker. A second support member has an inner edge that is secured to a lower portion of the voice coil former and an outer edge that is secured to the frame of the loudspeaker. In this type of loudspeaker, it is often desirable to reduce the outer-most dimensions of the frame of the loudspeaker to facilitate mounting the loudspeaker in various environments where space may be limited."}
{"text": "The present invention is a weatherproof illuminated tape that may be used for decorative lighting or other purposes.\nA common form of decorative lighting, often referred to as Christmas lights, employs a string of lights with individual bulbs or mini-lights mounted at spaced intervals in individual sockets.\nThe attraction of this type of lighting lies in the fact that it is flexible and can be mounted to surfaces of practically any shape, e.g. on Christmas trees or wreaths or around windows or doors. But, such lights, if not used carefully, are a known fire hazard in certain applications, e.g. use on indoor Christmas trees. Also, special measures must be taken if the lights are to be used out of doors, i.e. special outdoor lights (generally more expensive) which are water resistant must be used.\nChao U.S. Pat. No. 3,894,225 proposes sandwiching an elongated string of mini-lights between a pair of flexible transparent strips of plastic tape. The bulbs are connected in parallel between positive and negative thin wire leads that extend longitudinally along edges of the tape. The bulbs project outwardly through individual holes formed in the upper tape strip.\nScheib U.S. Pat. No. 4,439,818 proposes a similar structure, only using LED's instead of mini lamps. Positive and negative wires are laid out on adhesive tape and LED's, which are connected across the leads, project outwardly through holes formed in the tape.\nThe Chao and Scheib structures have inherent drawbacks. Fabrication of either product requires that the lamps or LED's be mounted on the carrier wires so that, when the electrical strip is joined to the upper tape strip, the lamps or LED's are in exact registry with the holes in the tape. Moreover, unless the holes are the exact size of the lamps, portions of the wires will be left exposed, which may produce a risk of shock. To solve the alignment problem, each of the bulbs or LED's could be positioned in the tape holes and individually wired in place. However, manufacture of the strip in such a manner would be labor intensive and therefore impractical for commercial purposes. Alternatively, the holes could be made larger to relieve the problems of alignment; but this will increase the problems associated with exposed wiring. Another problem is the fact that the bulbs or LED's project and are therefore exposed to damage. For such reasons, while Chao and Scheib propose the idea of a strip, neither appears to disclose a practical structure for doing so on a mass production basis."}
{"text": "The invention relates to chairs and the like, and more specifically to chairs that have arms with an adjustable height."}
{"text": "(a) Field\nThe present disclosure relates to a driving test system for a moving object that detects motion characteristics of the moving object traveling along a preset route and determines how the moving object, such as a vehicle, is functioning and driving.\n(b) Description of the Related Art\nIn general, an autonomous test vehicle is a vehicle that is forced to drive continuously to evaluate the vehicle's driving performance so that intended test results can be obtained. This vehicle is used to perform a forced driving operation on a test road such as a Belgian road with a rough surface.\nThus, the autonomous test vehicle does not require a human driver to manually test-drive the vehicle. This saves the operator the trouble and risk of test driving, allows for severe driving tests, and improves the reliability of test results. Therefore, research, development, and studies on techniques and methods for automatically controlling testing and driving are ongoing.\nConventionally, the above autonomous test vehicle self-controls its driving by enabling a vision sensor unit installed at the front of the test vehicle to recognize road lanes and detect an approaching object.\nThat is, the driving of the autonomous test vehicle is controlled in response to a control tower's radio control signals, which are input through an antenna and radio transmitter/receiver installed on the vehicle, and a driving control unit automatically controls the driving of the vehicle through a pedal controller in response to an approaching object and lane recognition signals, which are input from an object detector and vision sensor unit using various sensors.\nHowever, the operator's personal subjective view may influence a vehicle driving test, the accuracy of vehicle driving may be reduced, and the cost of labor may be increased.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "Administration of an anticancer drug such as Iressa (generic name: gefitinib) or Tarceva (generic name: erlotinib) causes a serious side effect called diffuse alveolar damage in some cases. Also in acute exacerbation of idiopathic pulmonary fibrosis, the serious side effect called diffuse alveolar damage is caused in some cases. Diffuse alveolar damage shows resistance to treatment, has a recurrent nature, and exhibits repeated exacerbation during treatment, resulting in serious exacerbation at an early stage. Thus, diffuse alveolar damage is a symptom to which very careful attention should be paid.\nJapanese are reported to develop such diffuse alveolar damage at high frequency (Non-patent Document 1). For example, in cases of gefitinib, the frequency is 0.3% according to surveys in countries other than Japan, while the frequency is 3.98% in Japanese, which indicates a not less than 10 times higher frequency in Japanese. In cases of erlotinib, the frequency is 0.2% according to surveys on Asians, while the frequency is 2.7% according to a survey on Japanese, which again indicates a not less than 10 times higher frequency in Japanese. Further, Japanese patients with idiopathic pulmonary fibrosis are reported to show a higher rate of acute exacerbation and a higher case fatality rate than patients in other countries (Non-patent Document 2). It has been assumed that the number of deaths due to diffuse alveolar damage in Japan is as much as several thousand per year. It should be noted, of course, that death due to the side effect of drugs is a serious problem also in countries other than Japan, apart from its frequency."}
{"text": "1. Field of the Invention\nThe present invention relates to an operational amplifier, and more particularly, to an operational amplifier capable of driving large capacitive loads with low power consumption.\n2. Description of the Prior Art\nOperational amplifiers have been applied extensively in the field of electrical devices and electronics, such as inverter amplifiers, integrators, and filter circuits, to name just a few instances. With the rapid scaling in CMOS processes, supply voltages in VLSL have been dramatically reduced in recent years. Acting as a fundamental block in most analog systems, operational amplifiers are required to achieve high gain and large bandwidth simultaneously in low-voltage applications. Since conventional cascade amplifiers, which increase the gain by stacking up transistors, are not suitable in low-voltage design due to small voltage swings, more circuit designers are aware of the importance of multi-stage amplifiers, which boost the gain by increasing the number of gain stages horizontally. However, all multi-stage amplifiers suffer close-loop stability problems due to their multiple-pole nature in the small-signal transfer functions. Therefore, many frequency compensation topologies have been proposed to ensure the stability of the multi-stage amplifiers. Generally, the operational amplifier applied in the conventional driver chip is normally a two-stage amplifier having a first stage amplifying circuit for gain enhancement and a second stage output circuit for driving the capacitive or resistive load. However, three-stage operational amplifiers are also gaining popularity.\nThe most relevant characteristics of an amplifier circuit are usually gain and bandwidth. There is an inverse relationship between the gain and the bandwidth of amplifiers. In general, higher gain values are associated with lower bandwidths, and lower gain values are associated with higher bandwidth. The performance of an operational amplifier is characterized by its transfer function which can be obtained by applying small-signal analysis. Reference is made to FIG. 1 for a transfer function of an exemplary two-stage amplifier. There is a relatively constant gain from DC to a frequency of the first dominant pole ωP1. When the frequency rises above ωP1, the gain begins to fall sharply. The maximum available bandwidth is related to the second non-dominant pole ωP2. It may be desirable to adjust the frequency of poles ωP1 and ωP2 for different applications. Various compensation techniques, such as Miller compensation or Ahuja compensation, are known for adjusting the frequency of the poles of the amplifier. Miller compensation employs a feedback capacitor connected across an input and output of the second amplifier stage. In Ahuja compensation, a current gain device is added in a feedback loop of the second amplifier stage.\nReference is made to FIG. 2 for a block diagram of a prior art two-stage Miller Compensation (MC) amplifier 10. The two-stage MC amplifier 10 includes a first-stage amplifier 11, a second-stage amplifier 12, and a compensation capacitor Cm. The transconductance, output resistance, and lumped output parasitic resistance of the first gain stage are notated by gm1, ro1, and Cp1, and those of the second gain stage are notated by gm2, ro2, and CL. The compensation capacitor Cm is coupled between the input and the output ends of the second-stage amplifier 12. By introducing the compensation capacitor Cm, the capacitance of the second-stage amplifier 12 appears much larger from its input, thereby shifting the first dominant pole ωP1 to a lower frequency and the second non-dominant pole ωP2 to a higher frequency.\nHowever, the capacitor Cm functions as a short-circuited path at high frequencies, and the combination of the capacitor Cm and the second-stage amplifier 12 creates a diode-connected transistor. In this case, any noise from a reference potential is transferred to the second-stage amplifier 12. In addition, the MC amplifier 100 has a poor power supply rejection ratio (PSRR) during high frequency operations. Therefore, if a good PSRR is required, the two-stage MC amplifier 10 is insufficient for desirable operation.\nReference is made to FIG. 3 for a block diagram of a prior art three-stage Nested Miller Compensation (NMC) amplifier 20. The three-stage NMC amplifier 20 includes a first-stage amplifier 21, a second-stage amplifier 22, a third-stage amplifier 23, and compensation capacitors Cm1 and Cm2. The transconductance, output resistance, and lumped output parasitic resistance of the first gain stage are notated by gm1, ro1 and Cp1, those of the second gain stage are notated by gm2, ro2, and Cp2, and those of the third gain stage are notated by gm3, ro3, and CL. The compensation capacitor Cm1 is coupled between the input end of the second-stage amplifier 22 and the output end of the third-stage amplifier 23. The compensation capacitor Cm2 is coupled between the input end and the output end of the third-stage amplifier 23. Under the assumptions: (1) Cm1, Cm2 and CL>>CP1 and CP2 and (2) gm3>>gm1 and gm2, the three-stage NMC amplifier 20 is characterized by the small-signal transfer function ANMC(s) represented by:\n                                          A            NMC                    ⁡                      (            s            )                          =                              A                          D              ⁢                                                          ⁢              C                                ⁢                                    A                              D                ⁢                                                                  ⁢                C                                                                    [                                                      s                                          ω                                              p                        ⁢                                                                                                  ⁢                        1                                                                              +                  1                                ]                            ⁡                              [                                                                            s                      2                                        ⁢                                                                                            C                                                      m                            ⁢                                                                                                                  ⁢                            2                                                                          ⁢                                                  C                          L                                                                                                                      g                                                      m                            ⁢                                                                                                                  ⁢                            2                                                                          ⁢                                                  g                                                      m                            ⁢                                                                                                                  ⁢                            3                                                                                                                                +                                      s                    ⁢                                                                  C                                                  m                          ⁢                                                                                                          ⁢                          2                                                                                            g                                                  m                          ⁢                                                                                                          ⁢                          2                                                                                                      +                  1                                ]                                                                        (        1        )            \nwhere ADC is the DC gain equal to gm1gm2gm3ro1ro2ro3 \nand ωP1 is the dominant pole equal to (Cm1gm2gm3ro1ro2ro3)−1 \nTo stabilize the NMC amplifier 20, following dimension condition shall be obeyed:\n                              C                      m            ⁢                                                  ⁢            1                          =                  4          ⁢                      (                                          g                                  m                  ⁢                                                                          ⁢                  1                                                            g                                  m                  ⁢                                                                          ⁢                  3                                                      )                    ⁢                      C            L                                              (        2        )                                          C                      m            ⁢                                                  ⁢            2                          =                  2          ⁢                      (                                          g                                  m                  ⁢                                                                          ⁢                  2                                                            g                                  m                  ⁢                                                                          ⁢                  3                                                      )                    ⁢                      C            L                                              (        3        )            \nPole-splitting can also be achieved in the NMC amplifier 20 by introducing the compensation capacitors Cm1 and Cm2. However, the non-dominant pole depends on Cm2 and thus depends on the loading capacitance CL, as depicted in (1) and (3). When driving a large capacitive load, a large Cm2 is required, thereby shifting the non-dominant pole to a rather low frequency. Therefore, the bandwidth of the NMC amplifier 20 is poor. Moreover, the NMC amplifier 20 is not suitable for low-power design since the previously stated assumption gm3>>gm1 and gm2 may not be valid.\nAhuja frequency compensation scheme is another well-known frequency compensation for operational amplifiers. Reference is made to FIG. 4 for a block diagram of a prior art two-stage Ahuja compensation amplifier 30. The two-stage Ahuja compensation amplifier 30 includes a first-stage amplifier 31, a second-stage amplifier 32, a compensation capacitor Cm, and a current gain device Ig. The transconductance, output resistance, and lumped output parasitic resistance of the first gain stage are notated by gm1, ro1, and Cp1, and those of the second gain stage are notated by gm2, ro2, and CL. The two-stage Ahuja compensation amplifier 30 is characterized by the small-signal transfer function Aahuja(S) represented by\n            A      ahuja        ⁡          (      s      )        =            A              D        ⁢                                  ⁢        C              ⁢                  [                              s                          ω              z                                +          1                ]                              [                                    s                              ω                                  p                  ⁢                                                                          ⁢                  1                                                      +            1                    ]                [                                            s              2                                      ω              n              2                                +                      s                          [                                                ω                  n                                                  2                  ⁢                                                                          ⁢                  ξ                                            ]                                +          1                ]            \nwhere\nADC is the DC gain equal to gm1gm2ro1ro2 \nωP1 is the dominant pole equal to (Cmgm2ro1ro2)−1 \nωz is the non-dominant zero equal to\n      g          m      ⁢                          ⁢      c            C    m  \nζ is the damping factor equal to\n                              1          2                ⁢                                                                              C                  1                                ⁢                                  g                                      m                    ⁢                                                                                  ⁢                    c                                                                                                C                  L                                ⁢                                  g                                      m                    ⁢                                                                                  ⁢                    2                                                                                ⁡                      [                          1              +                                                C                  L                                                  C                  m                                                      ]                                              (        4        )            \nωn is the natural frequency equal to\n                                                        g                              m                ⁢                                                                  ⁢                c                                      ⁢                          g                              m                ⁢                                                                  ⁢                2                                                                        C              1                        ⁢                          C              L                                                          (        5        )            \nIn the two-stage Ahuja compensation amplifier 30, the compensation capacitor Cm and the current gain device Ig are coupled in series between the input and the output ends of the second-stage amplifier 32. By introducing the current gain device Ig, the second-stage amplifier 32 no longer becomes a diode-connected transistor while it is operating at high frequencies. This two-stage Ahuja compensation amplifier 30 can thus achieve a good PSRR at high frequencies. However, the two-stage Ahuja compensation amplifier 30 fails to achieve good compensation in certain cases."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nIt is commonplace in arthroscopic procedures to employ sutures and anchors to secure soft tissues to bone. Despite their widespread use, several improvements in the use of sutures and suture anchors may be made. For example, the procedure of tying knots may be very time consuming, thereby increasing the cost of the procedure and limiting the capacity of the surgeon. Furthermore, the strength of the repair may be limited by the strength of the knot. This latter drawback may be of particular significance if the knot is tied improperly as the strength of the knot in such situations may be significantly lower than the tensile strength of the suture material.\nTo improve on these uses, sutures having a single preformed loop have been provided. FIG. 1 represents a prior art suture construction. As shown, one end of the suture is passed through a passage defined in the suture itself. The application of tension to the ends of the suture pulls a portion of the suture through the passage, causing a loop formed in the suture to close. Relaxation of the system, however may allow a portion of the suture to translate back through the passage, thus relieving the desired tension.\nIt is an object of the present teachings to provide an alternative device for anchoring sutures to bone and soft tissue. The device, which is relatively simple in design and structure, is highly effective for its intended purpose."}
{"text": "Commercial VI improvers with dispersant properties are typically made with a particular amine that is selected to disperse soot, or with a particular amine that is selected to disperse sludge and varnish. However, unless all deleterious materials in an engine oil are effectively dispersed, oil filter plugging may occur. Combining two different amine dispersants may not always be effective for dispersing all of the deleterious materials that may cause filter plugging.\nAdditionally, certain commercial VI improvers with dispersant properties may be highly effective for dispersing soot at the levels generated in today's modem engines in Group I and Group II base oils. However, the treat rate for such VI improvers in high viscosity index base stocks, i.e., base stocks having a viscosity index of greater than 120, is too low to provide effective dispersant performance and relatively high shear stability.\nAccordingly, there continues to be a need for lubricant formulations and additives that provide wider dispersant properties for modern engines and are effective at relatively lower treat rates than commercial VI improvers having dispersant properties.\nIn accordance with exemplary embodiments, the disclosure provides a lubricating oil composition and methods of operating an internal combustion engine to provide improved engine operational performance. The lubricating composition includes a major amount of oil of lubricating viscosity; and a minor amount of at least one highly grafted, multi-functional olefin copolymer made by reacting an acylating agent with an olefin copolymer having a number average molecular weight greater than about 10,000 up to about 50,000 in the presence of a free radical initiator to provide an acylated olefin copolymer having a degree of grafting (DOG) of the acylating agent on the olefin copolymer of from above about 1.5 to about 3.0 wt. %. The acylated olefin copolymer is then reacted with an amine to provide the highly grafted, multi-functional olefin copolymer.\nIn another exemplary embodiment, the disclosure provides a method for extending an oil drain interval for an internal combustion engine. The method includes lubricating the engine with a lubricant composition including a major amount of oil of lubricating viscosity, and a minor amount of at least one highly grafted, multi-functional olefin copolymer made by reacting an acylating agent with an olefin copolymer having a number average molecular weight greater than about 10,000 and less than about 50,000 in the present of a free radical initiator to provide an acylated olefin copolymer having a degree of grafting (DOG) of the acylating agent on the olefin copolymer ranging from about 1.5 to about 3.0 wt. %. The acylated olefin copolymer is then reacted with an amine to provide the highly grafted, multi-functional olefin copolymer, and the engine is operated using the lubricant composition.\nYet another exemplary embodiment of the disclosure provides a method of reducing oil filter plugging for an oil filter of an internal combustion engine. The method includes lubricating the engine with a lubricant composition including a major amount of oil of lubricating viscosity, and a minor amount of at least one highly grafted, multi-functional olefin copolymer made by reacting an acylating agent with an olefin copolymer having a number average molecular weight greater than about 10,000 and less than about 50,000 in the present of a free radical initiator to provide an acylated olefin copolymer having a degree of grafting (DOG) of the acylating agent on the olefin copolymer ranging from about 1.5 to about 3.0 wt. %. The acylated olefin copolymer is then reacted with an amine to provide the highly grafted, multi-functional olefin copolymer, and the engine is operated using the lubricant composition.\nAccordingly, a primary advantage of the exemplary embodiments may be an increased in oil change intervals due to a lower viscosity increase in the lubricating oil for an engine operating with the lubricating oil composition. Another advantage is that a copolymer dispersant VI improver made with a certain mixture of amines may provide synergistic performance for handling a wider range of deleterious materials that may accumulate in a lubricant during engine operation."}
{"text": "A wheel of a motor vehicle typically has a wheel hub that is rotatably supported on a wheel carrier (e.g. an axle arm or an axle funnel). A tire pressure regulation apparatus is provided for some applications, for example for heavy commercial vehicles, in order to set the air pressure in the tire mounted on the wheel hub to a desired value. In order to be able to supply the tire with the compressed air required for this purpose, the compressed air has to be transferred, for example starting from a central compressed air source, from the respective wheel carrier to the wheel hub supported thereat. The same applies accordingly in the opposite direction for deflating the tire. This is achieved by the rotary feedthrough which comprises a substantially ring-shaped intermediate space between the wheel carrier and the wheel hub in which a compressed air annular space is formed. It communicates both with a compressed air passage provided in the wheel carrier and with a compressed air passage provided in the wheel hub. In this manner, the named compressed air passages are also connected to one another when the wheel hub is rotating so that compressed air can be supplied to the tire via the compressed air passages, independently of the rotational position of the wheel hub.\nVery high pressures are consequently generated for the desired tire pressure regulation, for example, up to approximately 7 bar. A compressed air annular space formed in the intermediate space between the wheel carrier and the wheel hub must therefore be reliably sealed with respect to the environment and in particular with respect to the remaining intermediate space. This can take place, for instance, by means of sealing rings, in particular by means of radial shaft sealing rings, which bound the compressed air annular space at both sides and in particular with respect to annular lubricant spaces which are formed in the remaining intermediate space and which can serve, for example, for lubricating the sealing rings and/or bearings for supporting the wheel hub at the wheel carrier.\nSo that the pressure in the tire of a respective motor vehicle wheel can be varied, controllable valves are employed at the wheel hub, that is, in particular valves that can be selectively controlled either to let compressed air pass through into a respective filling direction or deflating direction or to block the passage of compressed air. In this respect, systems for regulating tire pressure can, for example, be configured as single-passage systems, i.e. having only a single compressed air passage in the wheel hub which leads to the valve, or as multi-passage systems in particular two-passage systems.\nWith two-passage systems, a compressed air passage in the wheel hub additionally leading to a valve exclusively serves for controlling the switch position of the valve. Two-passage systems have the advantage that a standard valve can be used at the tire whose air pressure should be able to be regulated. As a rule, such a standard valve comprises a compressed air port, a working port (2/2-way valve) and optionally a venting connection (3/2-way valve). For controlling the valve, that is in particular for blocking or (directed) opening of the valve, the valve can also have a control port which is actuated via the additional compressed air passage. In this respect, the valve can be selectively controlled via the pressure applied to the control port or via the power hereby caused at the control port to connect the compressed air port or the venting connection to the working port for the passage of compressed air or to block these connections with respect to one another. In so doing, the valve is preferably preloaded in a blocking state in order to maintain the tire pressure on a failure of the pressure or of the power at the control port.\nThe disadvantage of such a two-passage system is that the rotary feedthrough must comprise at least two compressed air passages in the wheel hub so that at least two compressed air annular spaces must also be formed in the intermediate space between the wheel carrier and the wheel hub in order to connect the compressed air passages of the wheel hub to associated compressed air passages of the wheel carrier. These compressed air annular spaces, which can have completely different pressures, must be reliably sealed, in particular also with respect to one another; for instance, in that respective sealing rings are provided between them and at the two sides.\nTo minimize the wear of the sealing rings sliding along the wheel carrier or the wheel hub, lubricants, for example lubricating grease or lubricating oil, can be guided through lubricant passages from the wheel carrier into annular lubricant spaces which are provided in the named intermediate space between the wheel carrier and the wheel hub in addition to the compressed air annular spaces. The lubricant can hereby move to the respective sealing rings bounding the respective annular lubricant space. The wear can admittedly thereby be reduced, but not the comparatively high number of wear parts. Such an arrangement also requires a relatively large construction depth (extent of the named intermediate space along the direction of arrangement of the sealing rings).\nA rotary feedthrough of such a two-passage system is known from the laying-open document DE 10 2009 057 158 A1.\nIn contrast, the valve position in single-passage systems is controlled by pressure impulses within the compressed air passage of the wheel hub. Single-passage systems can thus admittedly have a smaller construction depth of the named intermediate space and fewer wear parts since only a single compressed air passage is provided in the wheel hub for supplying the tire with compressed air and only a single compressed air annular space therefore also has to be formed in the intermediate space between the wheel carrier and the wheel hub. However, a single-passage system requires the use of a special valve which can only be controlled by means of the compressed air port which thus simultaneously functions as a compressed air inlet and a control port. However, such a special valve causes higher costs and necessitates a more complex control of the air pressure in the compressed air passage for a defined generation of the control pressure bursts. Furthermore, there is the risk of a certain proneness to disturbance."}
{"text": "Planarization of semiconductor substrates is becoming more important as the number of layers used to form a semiconductor device increases. Nonplanar semiconductor substrates have many problems including difficulty in patterning a photoresist layer, formation of a void within a film during the film deposition, and incomplete removal of a layer during an etch process leaving residual portions of the layer, which are sometimes called \"stringers.\" A number of planarization processes have been developed and include chemical-mechanical polishing.\nFIGS. 1 and 2 include illustrations of a part of one type of a chemical-mechanical polisher that is used to polish semiconductor substrates. FIG. 1 is a cross-sectional view of a chemical-mechanical polisher 10. The polisher 10 has a platen 14 and a polishing pad 11 attached to the platen 14 with an adhesive compound (not shown). Above the polishing pad 11 are substrate holders 12, and each substrate holder 12 has a semiconductor substrate 13. The polisher 10 also includes a polishing slurry and a slurry feed, both of which are not shown. The polishing pad 11 may be made of a porous polyurethane material that has a relatively uniform thickness of about 1-2 millimeters. FIG. 2 includes a top view illustrating the relationships of motion between the polishing pad 11 and the substrates 13. During polishing, the polishing pad 11 rotates counterclockwise or clockwise, but the substrates 13 typically rotate in the same direction as the polishing pad 11. While the substrates 13 and polishing pad 11 are rotating, the substrates 13 are being oscillated back and forth across the polishing pad. The oscillating motion covers a distance called an oscillating range and is performed at an oscillating velocity. While the polishing is being performed, the polishing slurry may be recycled.\nIn actual use, chemical-mechanical polishing typically has nonuniform polishing rates across a substrate surface. In many cases, the polishing rate near the edge of the substrate is higher than the polishing rate near the center of the substrate because the relative velocity between polishing pad and the substrate is higher near the edge of the substrate compared to the center of the substrate. Therefore, some area of the substrate near the center may be underpolished, some area of the substrate near the edge may be overpolished, or both.\nThe polishing pad may contribute to the nonuniformity. A brief overview of the formation of polyurethane polishing pads is now presented. Polyurethane polishing pads are typically formed by reacting the chemicals that form polyurethane within a cylindrical container. After forming a cylindrical-shaped piece of polyurethane, the piece is cut into slices that are subsequently used as polishing pad. The polishing pad typically has pores that have a size of about 100-200 microns. Although the pores may vary in size, the average pore size for any region of the polishing pad is typically about the same as any other region of the polishing pad. As used hereinafter, this type of prior art polishing pad is referred to as a conventional polishing pad. The nonuniformity occurs because the edge of the substrate is moving faster relative to the polishing pad compared to the center of the substrate and the conventional polishing pad does not have a feature to compensate of the polishing nonuniformity.\nThe prior art has addressed the problem of nonuniformity polishing by modifying a conventional polishing pad by forming a pattern within the polishing pad. These polishing pads include forming a variety of geometric patterns including openings. It should be kept in mind that polishing pads are typically porous, and the pores are formed during the reaction to form the polishing pad material. As used in this specification, openings are distinguished from pores because openings are formed within the pad after the reaction to form the polishing pad material has occurred. A conventional polishing pad has pores but does not have any openings. The prior art polishing pad with openings typically have a width on the order of centimeters, or the prior art polishing pad has a density of openings that decreases with the distance from edge of the polishing pad."}
{"text": "Conventional printing devices, e.g., inkjet and laser printers, photocopiers, facsimile machines, etc., typically implement a best-fit four color printing process in which cyan, magenta, yellow and black colored recording materials, e.g., ink, toner, etc., are combined in an attempt to obtain a desired color. It has been found that the best-fit four color printing process, while suitable for a great number of printing applications, may be unable to produce certain colors with suitable accuracy.\nAs an example, when an entity, e.g., person, business, corporation, etc., possesses a logo that includes a color that the best-fit color matching process is unable to match in a substantially adequate manner, the only option may be for the entity to rely on an outside printing vendor to print the logo in the desired color.\nReliance upon an outside printing vendor to print the logo on various media, e.g., letterhead, business cards, etc., may result in certain drawbacks and disadvantages. In one respect, the entity may need to order substantially large quantities of the various media to ensure that an adequate supply is readily available. One result of maintaining substantially large quantities of the various media is that large portions of the supply may be discarded, for example, if the letterhead changes, if an employee leaves the company, etc. In addition, the costs involved in purchasing the various media from a printing vendor may be relatively high.\nIt is generally known to include e-labels on cartridges, e.g., inkjet, toner, etc. E-labels may generally be defined as a memory storage element, e.g., EEPROM, PROM, radio frequency identification tags, etc., and may contain information related to the identity of the ink contained in the cartridge, e.g., ink type, color, date of manufacture, and the like. The information may be retrieved by, for example, a printing device controller."}
{"text": "The present invention relates to filters, and particularly to underdrains in gravity filters. More particularly, the present invention relates to underdrains in gravity filters having a cap or plate supporting a filter media.\nGravity filters frequently employ underdrain blocks to channel fluid away from the filter bed and to direct backwash fluids evenly into the filter bed during a backwash cycle. Underdrain blocks are assembled end to end forming rows, also known as laterals. Several laterals are placed next to one another to define the filter bottom. A filter media (e.g., granular activated carbon, anthracite, coal, magnesium oxide, ilmenite, sand including garnet, silica or quartz, etc.) placed on top of the laterals completes the filter bed and performs a majority of the filtration.\nIn operation, the filter bed receives a flow of fluid to be filtered and passes it through the filter media to the underdrain blocks. The size and arrangement of the filter media allows it to filter out the undesired contaminates, with finer filter media being able to filter out smaller contaminates. The flow, or filtrate, enters the underdrain blocks through openings generally provided in the tops of the underdrain blocks. The underdrain block laterals provide a convenient channel or flow path, allowing removal of the filtrate from the filter bed.\nPeriodically, a backwash cycle is initiated to clean the filter media. During a backwash cycle, fluid flow reverses, thus flowing into the underdrain blocks and upwardly through the filter media to loosen and remove contaminates.\nIt is desirable to have large openings in the filter blocks to allow free filtrate flow into the blocks. However, large openings are incapable of supporting the filter media and instead allow the media to pass into the underdrain blocks, which is undesirable. Small openings on the other hand, are capable of supporting the filter media, but are easily plugged by the filter media disposed above and upstream of the underdrain blocks. Also, small openings provide greater restriction to flow than large openings even when unclogged. Therefore, it is desirable to provide an underdrain block capable of supporting the filter media above, while simultaneously reducing the likelihood of media entering the underdrain block, and minimizing plugging of the openings. It is also desirable to provide an underdrain block having a high percentage of open surface area, thereby allowing a free fluid flow therethrough.\nThus, according to the present invention, a filtration system provides a block including a first wall and a second wall substantially parallel to the first wall. A top surface that is substantially perpendicular to the first wall is disposed between, and coupled to, the first and second walls to at least partially define an inner chamber. The top surface further includes an orifice defining a passageway between a block exterior and the inner chamber. A plate including a plurality of slots, each slot having a width of less than about 0.065 inches, and a slot outlet, is coupled to the top surface of the block. The slot outlets in the plate are substantially coplanar and the plate, when coupled to the block, substantially covers the orifice.\nThe invention further provides a method of manufacturing the filtration system. The filtration system includes an underdrain block, a retaining plate, and a filter media. The underdrain block includes an outer wall having a first side and a second side substantially parallel to the first side. The outer wall also includes a top surface substantially perpendicular to the first and second sides. A plurality of orifices are defined within the top surface. The retaining plate has an upper surface and a lower surface. The method includes the steps of forming the retaining plate such that the upper surface and the lower surface are substantially parallel to one another. The method further includes laser etching the plurality of slots into the retaining plate such that the area of the slots is at least equal to the area of the orifices, and such that the slot outlets are coplanar. The method further comprises aligning the retaining plate with the underdrain block, and coupling the retaining plate to the underdrain block such that the retaining plate at least partially covers the plurality of orifices of the underdrain block.\nIn preferred embodiments, the retaining plate is manufactured from a stainless steel plate that is less than one-eight of an inch thick. The plurality of slots are preferably laser-etched through the plate and have a width between 0.002 inches and 0.025 inches. A series of longitudinal and cross members defined by the block provide support to the plate when attached to the block. In preferred embodiments, a plurality of screws attach the retaining plate to the block.\nAdditional features and advantages will become apparent to those skilled in the art upon consideration of the following detailed description of preferred embodiments exemplifying the best mode of carrying out the invention as presently perceived."}
{"text": "The present invention relates generally to interactive electronic toys, and more particularly to a uniquely configured sensor and associated electronic circuitry which may be incorporated into interactive electronic toys and games (including dolls and remote controllers such as joysticks) and is operative to produce various visual and/or audible outputs or signal transmissions corresponding to the level/position of the toy relative to a prescribed plane.\nThere is currently known in the prior art a multitude of interactive electronic toys which are capable of producing a wide variety of visual and/or audible outputs. In the prior art toys, these outputs are typically triggered as a result of the user (e.g., a child) actuating one or more switches of the toy. The switch(es) of the prior art toys are most typically actuated by pressing one or more buttons on the toy, opening and/or closing a door or a hatch, turning a knob or handle, inserting an object into a complimentary receptacle, etc. In certain prior art interactive electronic toys, the actuation of the switch is facilitated by a specific type of movement of the toy. However, in those prior art electronic toys including a motion actuated switch, such switch is typically capable of generating only a single output signal as a result of the movement of the toy.\nThe present invention provides a uniquely configured sensor and associated electronic circuitry which is particularly suited for use in interactive electronic toys and games, including dolls and remote controllers such as joysticks. The present sensor is specifically configured to generate a multiplicity of different output signals which are a function of (i.e., correspond to) the level/position of the toy relative to a prescribed plane. Thus, interactive electronic toys and games incorporating the sensor and associated electronic circuitry of the present invention are far superior to those known in the prior art since a wide variety of differing visual and/or audible outputs and/or various signal transmissions may be produced simply by varying or altering the level/position of the toy relative to a prescribed plane. For example, the incorporation of the sensor and electronic circuitry of the present invention into an interactive electronic toy such as a spaceship allows for the production of differing visual and/or audible outputs as a result of the spaceship being tilted in a nose-up direction, tilted in a nose-down direction, banked to the left, and banked to the right. As indicated above, the output signals generated by the sensor differ according to the level/position of the sensor relative to a prescribed plane, with the associated electronic circuitry of the present invention being operative to facilitate the production of various visual and/or audible outputs corresponding to the particular output signals generated by the sensor.\nIf incorporated into a joystick or other remote controller, the present sensor and associated electronic circuitry may be configured to facilitate the production of the aforementioned visual and/or audible outputs, and/or generate electrical/electronic signals, radio signals, infrared signals, microwave signals, or combinations thereof which may be transmitted to another device to facilitate the control and operation thereof in a desired manner. The frequency and/or coding of the radio, microwave, or electrical/electronic signals and the coding of the infrared signals transmitted from the joystick or other remote controller would be variable depending upon the level or position of the same relative to a prescribed plane. Moreover, the present electronic circuitry may be specifically programmed to memorize or recognize a prescribed sequence of movements of the sensor relative to a prescribed plane. More particularly, a prescribed sequence of states or output signals generated by the sensor corresponding to a prescribed sequence of movements thereof, when transmitted to the electronic circuitry, may be used to access a memory location in the electronic circuitry in a manner triggering or implementing one or more pre-programmed visual and/or audible functions or effects and/or the transmission of various electrical (hard wired), infrared, radio, or microwave signals to another device for communication and/or activation of various functions thereof. These, and other unique attributes of the present invention, will be discussed in more detail below."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor optical device having an active layer of a quantum well structure, such as an oscillation polarization mode selective semiconductor laser whose oscillation polarization mode can be changed by a direct modulation method with its dynamic wavelength fluctuation being oppressed even at the time of high-speed modulation and an polarization-insensitive optical amplifier which can amplify light substantially irrespective of its polarized condition, a use method of the semiconductor optical device and a method for driving the semiconductor optical device. The present invention also relates to a light source apparatus, an optical communication method and an optical communication system using the optical semiconductor device with an active layer of a quantum well structure.\n2. Related Background Art\nAs an oscillation polarization-mode selective dynamic single-mode semiconductor laser, the following device has been developed and proposed. The oscillation polarization mode selective device has a structure that can be modulated by a digital signal which is produced by superposing a small-amplitude digital signal on a bias injection current (see, for example, Japanese Patent Laid-Open Application No. 7(1995)-162088). The device is a distributed feedback (DFB) laser in which a distributed reflector of a grating is provided in a semiconductor laser resonator or cavity and wavelength selectivity of the grating is utilized. In the device, a bulk active layer is used, or strain is introduced into an active layer of a quantum well structure and its Bragg wavelength is located at a position shorter than a peak wavelength of a gain spectrum, so that gains for transverse electric (TE) mode and transverse magnetic (TM) mode are approximately equal to each other for light at wavelengths close to its oscillation wavelength, under a current injection condition near its oscillation threshold. Further, a plurality of electrodes are arranged and currents are unevenly injected through those electrodes.\nAn equivalent refractive index of the cavity is unevenly distributed by the uneven current injection, and its oscillation occurs in one of the TE mode and the TM mode and at a wavelength which satisfies a phase matching condition and takes a minimum threshold gain. When the balance of the uneven current injection is slightly changed to vary a competitive relation of the phase matching condition between the TE mode and the TM mode (i.e., in which mode is the threshold lower than the other's threshold in a state which satisfies the phase matching condition), the oscillation polarization mode and wavelength of the device can be switched.\nIn that semiconductor device, an antireflection coating is provided on one end facet to asymmetrically employ effects of the uneven current injection into its output-side portion (a modulation current is not injected into the output-side portion with one electrode so as not to fluctuate an output power) and its modulation-electrode portion (a portion where the other electrode is arranged). Alternatively, lengths of those electrodes are made different from each other to obtain a structural asymmetry.\nFurther, Japanese Patent Laid-Open Application No. 2-117190 discloses a semiconductor laser apparatus in which two semiconductor structures are arranged serially or in parallel. One of the semiconductor structures principally oscillates or amplifies a light wave in a predetermined polarization mode, and the other one chiefly oscillates or amplifies a light wave in another polarization mode. Those semiconductor structures are formed on a common layer plane or in parallel layer planes.\nIn the polarization selective DFB laser in which the oscillation polarization mode is selected depending on its above-discussed phase condition (i.e., the oscillation in a polarization mode having a lower threshold is selected at a wavelength which satisfies its phase condition), it is important to approximately equalize gains for the TE mode and the TM mode with each other in its cavity. As an active layer for such a purpose, there have been proposed a method of forming a bulk active layer with about the same TE-mode and TM-mode gains, a method of arranging an active layer with a dominant TM-mode gain obtained by the introduction of strain or the like thereinto in sequence with an active layer for generating a TE-mode gain (see Japanese Patent Laid-Open Application No. 2(1990)-117190).\nThose structures, however, have the following drawbacks. In the bulk active layer, though the TE gain can be made equal to the TM gain, its threshold current tends to increase. Therefore, its threshold current needs to be reduced by its quantization or the like.\nOn the other hand, in the method of Japanese Patent Laid-Open Application No. 2-117190 in which active layers with dominant TE and TM gains are serially or in parallel are arranged, the number of its growth steps is increased and its fabrication process is complicated since active layers with different strains must be formed.\nFurthermore, in those active layers, a wavelength band width, over which gains are substantially constant or uniform under a near-threshold condition, is narrow. For example, in a simple quantum well structure, a wavelength range, over which the gain is within a three-dB-down value from its peak and which includes a peak wavelength, is approximately 20 nm. Therefore, considering into account difficulties of the wavelength control near the gain peak wavelength of the active layer and the control of a grating pitch, it is hard to fabricate a low-threshold current laser with a good reproducibility and a stable oscillation wavelength by precisely controlling those gain peak wavelength and grating pitch.\nIn addition, when considering it necessary that respective wavelength ranges, over which TE-mode and TM-mode gains are substantially uniform and balance with each other, overlap on each other, it is desirable that a wavelength range, over which a gain for each polarization mode is approximately constant, is sufficiently wide. Furthermore, when an array laser consisting of a plurality of lasers with different oscillation wavelengths are to be achieved by using a common active layer, it is necessary that a wavelength range, over which each gain is approximately uniform, be expanded."}
{"text": "Many types of conveyors are widely known in the art. Conveyors typically comprise a fixed or articulated frame member, or boom, with an endless conveyor belt or chain located about the boom. Conveyors are generally used to transport materials, such as bundles of shingles and other roofing materials from, i.e., a truck to a roof top. Another form of conveyor system that may be utilized in conjunction with the present invention is a latter truck, such as the latter trucks used by fire departments.\nOne inherent problem with conveyors involves the angle of the conveyor in relation to an upper surface, such as a roof. This problem is especially evident when transporting, e.g., shingles onto a residential housing roof or other pitched surface. Since the conveyor is angled upwardly when transporting materials, the end of the conveyor may be vertically-spaced from the roof surface. As such, the person unloading the conveyor may have to reach above his or her head to receive the materials carried by the conveyor. Additionally, pitched roofs render a person's footing unsafe, especially when removing heavy items from the conveyor.\nThere is a need for a conveyor that provides a more level and safer approach angle for the person unloading materials from the conveyor. Accordingly, the present invention is hereby submitted."}
{"text": "Swimming pool floats come in a variety of sizes, shapes and materials. Regardless of this variety, the floats are all soft, light in weight and in need of drying after use.\nOne type of pool float is inflated with air. This type of float includes mattresses, chairs, etc. Another type of pool float is made of foam. The higher quality foam floats have a rubber or plastic coating over the foam.\nEquipment owners typically store the floats in a dry place, out of the pool, between uses. This extends the useable life of the floats, which can be shortened due to exposure to pool chemicals and sunlight. Most equipment owners merely stack the floats one on top of another. Foam floats stowed in this manner bend or become indented. The bends and indentations may remain in the float for long periods of time when the float is removed from storage.\nAlso, stacking the floats may not allow the floats to adequately dry, leading to mildew problems.\nThus, what is needed is a low cost rack for storing pool floats when the floats are not in use."}
{"text": "The goal of crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits may include greater yield, resistance to insects or pests, tolerance to heat and drought, better agronomic quality, higher nutritional value, growth rate and fruit or pod properties.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant or plant variety. A plant cross-pollinates if pollen comes to it from a flower of a different plant variety.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants of different varieties produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform varieties requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods that have been used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more plants or various other broad-based sources into breeding pools from which new lines are developed by selfing and selection of desired phenotypes. The new lines are evaluated to determine which of those have commercial potential.\nOne crop species which has been subject to such breeding programs and is of particular value is garden bean (Phaseolus vulgaris (snap)). Beans are annual, warm-season legumes. Garden beans, also known as green beans, snap beans, or pole beans, are grown primarily for their pods, which are harvested for consumption in their succulent form, whereas dry beans (Phaseolus vulgaris (dry)), lima beans (Phaseolus limensis), and soybeans (Glycine max) are usually grown for the seed itself. In addition, the bean leaf is occasionally used as a leaf vegetable, and the straw is used for fodder.\nGarden beans are available in bush and pole varieties. Bush varieties form erect bushes 20-60 centimeters tall, while pole varieties form vines 2-3 meters long. The pods are typically 8-20 centimeters long, 1-1.5 centimeters wide, and either green, yellow, black or purple in color. Each pod generally contains 4-6 beans."}
{"text": "A context based arithmetic encoding CAE for binary shape coding of I (intra-frame) VOPs, P (unilaterally motion-compensated-predicted) VOPs and B (bidirectionally interpolated) VOPs (VOP=video object plane) is used in MPEG verification model version 8.0 (VM8.0), ISO/IEC JTC1/SC29/WG11, MPEG97/N1796.\nFIG. 1 shows an example of a GOF (group of frames) structure in MPEG-1 standard for an image sequence composed of I-, P- and B-VOPs. VOPs coded exclusively by the intra-frame technique are labeled as I, unilaterally motion-compensated predicted VOPs are labeled as P, and the bidirectionally interpolated VOPs between them are labeled as B. The example according to FIG. 1 shows four P-VOPs per GOF and two B-VOPs between I- and/or P-VOPs. The length of the GOF is not fixed, but instead is determined by the number of different image types. One essential feature of the CAE algorithm is that a chronologically following VOP (backward reference VOP) is always used as the reference value in inter-CAE shape coding of a current B-VOP, namely for both the binary shape reference and the shape mode reference. However, since B-VOPs of a random shape sometimes access only image contents of a chronologically preceding VOP as the reference (forward reference VOP), the binary shape of this B-VOP must be coded in the same manner as an I-VOP when using the current CAE algorithm. This leads to a considerable degradation of coding efficiency with such B-VOPs.\nAccording to the present invention, coding efficiency can be improved in particular for shape coding in bidirectionally interpolated VOPs (B-VOPs). The method according to the present invention provides time-adaptive shape coding for B-VOPs. During a context-based shape coding, a B-VOP can access either a chronologically preceding I-VOP or P-VOP (forward reference VOP) or a chronologically following I-VOP or P-VOP (backward reference VOP) as the reference value for shape coding.\nIn the method according to the present invention no changes in syntax are necessary and thus the method can easily be used with existing standards such as the MPEG-4 standard."}
{"text": "At fast food restaurants and at take-out counters, french fries are often sold in paperboard cartons. Many people like to apply ketchup to french fries. If it is applied directly to the fries in the container, they will tend to become soggy. Alternatively, one can obtain a small separate container of ketchup, in which a french fry may be dipped immediately before it is eaten. But this latter gourmet preference has the disadvantage of requiring two separate containers. With the human limitation of two hands, there is one hand to hold the container of fries and one to move the fries from container to mouth. So a stable and convenient surface is required to support the ketchup. But, since fast food customers are often on the move, they want to carry the fries with them as they walk or ride. A solution to the problem is a container for the fries, which includes an integral, discrete compartment for the ketchup."}
{"text": "1. Field of the invention\nThe present invention relates to rotary magnetic heads utilized in video tape recorders and digital audio recorders, and particularly relates to a method of inspecting the magnetic heads to find whether each of them are properly assembled onto a rotary drum or not.\n2. Related art\nIt is well known that what is called \"a double azimuthal head\" as shown in FIG. 3(a) is utilized for the rotary drum X of the video tape recorder or the digital audio tape recorder which records and reproduces at a reference speed and at the twice-faster speed, and two pairs of magnetic heads are assembled on the rotary drum in such manner that each head confronts the other in the diametric direction. In FIG. 3(b), a magnetic head A and a magnetic head C are the recording and reproducing heads for the first channel and a magnetic head B and a magnetic head D are the recording and reproducing heads for the second channel.\nIn a conventional inspection process of the rotary magnetic heads, each of the magnetic heads A through D is judged to be proper or improper by such method as shown in FIG. 5, where a magnetic tape is fed at the reference speed (Step S.sub.a) and, after a reference signal is recorded on the aforenamed magnetic tape by the magnetic head A (Step S.sub.b), the inspection apparatus is set to the reproducing mode to reproduce the reference signal which has been recorded by the aforenamed magnetic head A (Step S.sub.c) to judge whether the recorded signal is higher than the reference level (Step S.sub.d); if it is lower than the reference level, a troubleshooting measure, such as illumination of a warning lamp, is taken (Step S.sub.e).\nIn the aforenamed way, the aforenamed inspection process is repeated for the magnetic heads B through D till all the magnetic heads are checked whether they meet the reference value.\nAs FIG. 4 shows, each of the magnetic heads A through D has an azimuthal angle +.theta..sub.1, .theta..sub.2, -.theta..sub.2 or -.theta..sub.1 and, if these magnetic heads A through D are assembled onto a head base 1 shown in FIG. 3(a) in a incorrect way, the expected functions cannot be exercised.\nWhat has been inconvenient with the conventional inspection method as described above is that the magnetic heads A through D are judged only to be good or defective from their respective performance; whether or not the magnetic heads A through D are properly assembled onto the head base 1 is not inspected."}
{"text": "Vehicles, such as light-duty cars and trucks and heavy-duty tractor/trailers, may include “telematics” systems that monitor information describing the vehicle's location and diagnostic condition. Such telematics systems may include a conventional global positioning system (“GPS”) that receives signals from orbiting satellites and a processor that analyzes these signals to calculate a GPS “fix.” The fix describes the vehicle's location with an accuracy of about 10 meters or better and may include data such as the vehicle's latitude, longitude, altitude, heading, and velocity.\nTelematics systems may also include circuitry that monitors the host vehicle's diagnostic system. For example, light-duty automobiles and trucks beginning with model year 1996 may include an on-board diagnostic (“OBD-II”) system as mandated by the Environmental Protection Agency (“EPA”). OBD-II systems may operate under one of the following communication protocols: J1850 VPW (Ford); J1850 VPWM (General Motors); ISO 9141-2 (most Japanese and European vehicles); Keyword 2000 (some Mercedes and Hyundai vehicles); and CAN (a newer protocol used by many vehicles manufactured after 2004).\nOBD-II systems monitor the vehicle's electrical, mechanical and emissions systems and generate data such as vehicle speed (“VSS”), engine speed (“RPM”), engine load (“LOAD”), and mass air flow (“MAF”). This data may be processed by a vehicle's engine control unit (“ECU”) to detect malfunctions or deterioration in vehicle performance. In processing such data, the ECU may generate diagnostic trouble codes (“DTCs”), which are 5-digit codes (e.g., “P0001”) indicating electrical or mechanical problems with the vehicle. Most vehicles manufactured after 1996 make the information available through a standardized, serial 16-cavity connector located in the vehicle and referred to herein as an “OBD-II connector.” The OBD-II connector provides a serial interface to the vehicle's ECU and passes data and diagnostic codes to diagnostic equipment located external to the vehicle.\nLikewise, heavy-duty trucks may also include a diagnostic system referred to herein as a “truck diagnostic system,” which is analogous to the OBD-II systems present in light-duty vehicles. Truck diagnostic systems may operate under a communication protocol such as J1708/J1587 or J1939 and may include a connector located in the truck's interior and referred to herein as the “truck diagnostic connector.” The truck diagnostic connector, like the OBD-II connector, makes the diagnostic information available to diagnostic equipment located external to the vehicle."}
{"text": "It is known in the art to use dental appliances for the treatment of snoring and sleep apnea. For example and without being limitative, the dental appliances can be intraoral orthoses which prevent the mandible from moving posteriorly during sleep.\nOne type of known intraoral orthosis includes two dental trays (or bite forms) joined together so that the user's mandible projects forwardly from its normal position and thereby helps reducing snoring and sleep apnea. In such an orthosis, the dental trays are connected to one another through two lateral links which are configured to advance the user's mandible during a vertical movement between the maxilla and the mandible. In some embodiments, for connection to the dental trays corresponding to the maxilla or the mandible, the lateral links have a protrusion inserted in a through hole of the dental trays.\nOne of the common drawbacks of known intraoral orthoses is that the lateral links often disconnect from the dental tray to which they are engaged, when the intraoral orthosis is being manipulated. In such cases, the user must spend time to reconnect the disconnected lateral link to its respective dental tray.\nIn view of the above, there is a need for an improved mandibular protrusion device which, by virtue of its design and components, would be able to overcome or at least minimize some of the above-discussed prior art concerns."}
{"text": "The present invention relates generally to devices for recording data onto recordable optical discs, such as a CD-R (Compact Disc-Recordable) or DVD (Digital Video Disc-Recordable) where data are recorded on a packet-by-packet basis, and it relates more particularly to a packet-write-type optical disc recording device which records data in packets each having a predetermined data length.\nExamples of the conventionally-known methods for recording data on a CD-R include the \"disc-at-once\" method in which data for a whole optical disc are recorded in succession by single recording operation and the \"track-at-once\" method in which a plurality of sessions are allowed to be written onto a disc in an additive fashion (additive writing). All of these recording methods, however, involve recording of a great quantity of data per recording operation, consuming a managing memory area of 10 to 20 megabytes per recording operation. Thus, the known recording methods are not suited to applications where a relatively small quantity of data are written a plurality of times in an additive fashion.\nIn contrast, the so-called \"packet-write\" recording method, in which a small quantity of data are written onto an optical disc on a packet-by-packet basis, has the advantage that any necessary file can be additively written in much the same way as normally done on conventional hard discs, and hence it has been becoming increasingly popular in applications where a CD-R is used as an external storage device for a computer.\nFIG. 5 is a diagram illustrating a packet-write recording format. Each track is composed of a pre-gap, which is a managing area provided at the head of the track, and a plurality of packets following the pre-gap. In this format, data are written on the packet-by-packet basis; that is, each packet is written onto a disc, following an already-written packet. For this purpose, the pre-gap contains track descriptors, and a specific position to record a new packet is determined by reference to position information of a last packet written on the optical disc. In additively writing a new packet, it is impossible to physically link the new packet with the preceding already-written packet in a continuous manner, so that a pratically unreadable block called a \"link block\" is provided between every adjoining packets. Link information called a \"splice\" is stored in a data link period between the packets. The splice comprises information stored in a total of seven blocks: run-out information (Run-out) stored in two blocks at the end of a preceding packet; link information (Link) stored in one block at the head of a succeeding packet; and run-in information (Run-in) in four blocks of the succeeding packet.\nThe packet length is set by an external host computer. Because each packet comprises a plurality of blocks (sectors) each having a length of about 2K bytes, the packet length may be determined using 2K-byte units in principle However, due to the splice consuming a total of 14K bytes, the recording areas can not be used efficiently if the packet length is too short. For this reason, the packet length is normally set to 64K bytes, 128 bytes or the like so that the effective use rate of the recording areas exceeds 80 per cent.\nBecause of the nature of the packet-by-packet additive writing, the conventional optical disc recording devices conduct data recording onto an optical disc on the packet-by-packet basis. When one packet of data is to be recorded, a series of operations is generally required to store data for one packet from the host computer into an internal data buffer of a disc drive device, move a recording head to a particular writing position facing the disc (seek operation) and then initiate recording of the stored one-packet data by means of the recording head. The reason why the recording is initiated after storage of the one-packet data is that if there ocurs a break or interruption of data during recording, all data having so far been recorded would become errors.\nIf the size of a file is short enough to fit within a single packet, the above-mentioned recording method would present no significant problem; however, if the file is longer, a problem would arise that the seek operation frequently takes place, making recording of the file time-consuming. In addition, laser power of the recording head is frequently calibrated, which would result in shortage of a calibrating area. Such a problem may be lessened if the packet size is increased, but the increased packet size would considerably deteriorate the advantageous function of the packet-write recording method which allows a small quantity of data to be recorded in an additive fashion."}
{"text": "1. Field of the Invention\nThis invention generally relates to a cycle computer display apparatus for a bicycle. More specifically, the present invention relates to a cycle computer display apparatus having a touch screen.\n2. Background Information\nBicycling is becoming an increasingly more popular form of recreation as well as a means of transportation. Moreover, bicycling has become a very popular competitive sport for both amateurs and professionals. Whether the bicycle is used for recreation, transportation or competition, the bicycle industry is constantly improving the various components of the bicycle.\nRecently, bicycles have been equipped with electrical components to make riding easier and more enjoyable for the rider. Some bicycles are equipped with automatic shifting units that are automatically adjusted according to the riding conditions by a cycle computer or control unit. In particular, the front and rear derailleurs have recently been automated. Moreover, various electronic devices have been used to determine one or more operating parameters for controlling the derailleurs and providing information to the rider. Thus, the cycle computer or control unit of the bicycle not only provides information to the rider, but also controls various components of the bicycle.\nIn view of the above, it will be apparent to those skilled in the art from this disclosure that there exists a need for an improved cycle computer display apparatus. This invention addresses this need in the art as well as other needs, which will become apparent to those skilled in the art from this disclosure."}
{"text": "1. Field of the Invention\nThe present invention relates to a rotor including a magnetic holding structure and an electric motor including the rotor.\n2. Description of the Related Art\nA known synchronous electric motor includes a rotor having a plurality of magnets arranged on a surface thereof in a circumferential direction. When a higher speed of rotation is desired in this type of electric motor, it is necessary to provide sufficiently high holding strength for holding the magnets, in order to prevent the magnets from being detached from the rotor due to centrifugal force.\nJP-A-11-089142 discloses a rotor of an electric motor, which includes a rotational axis, a sleeve provided in an outer circumference of the rotational axis, a cylindrical permanent magnet provided on an outer circumference of the sleeve, and a holding member made of carbon-fiber-reinforced plastic and provided on an outer circumference of the permanent magnet, so as to cover the permanent magnet.\nJP-A-2013-169103 discloses a rotor of an electric motor, which includes a cylindrical protective tube provided so as to surround a plurality of magnets. The protective tube is attached to the rotor so as to have an increased diameter with filling pressure of resin, and therefore the protective tube is pressed against the magnets on a radially inside due to restorative force of the protective tube. According to this related art, the magnets can be held in position without an adhesive.\nIn order to rotate a rotor at a high speed, even greater holding strength on the magnets is necessary. For example, when adopting the holding member made of carbon-fiber-reinforced plastic as disclosed in JP-A-11-089142, the holding strength can be increased by thickening the holding member. However, since stress is generated concentrically on an inner diameter side, simply increasing the thickness of the holding member is not effective. In addition, when a thickened holding member is adopted, a distance between the magnets of the rotor and the stator (magnetic gap) becomes greater, resulting in the decreased torque. Further, it is expensive and a high level of skill is necessary to form a thickened fiber-reinforced plastic member, possibly causing defects in the internal structure and impairing the holding strength.\nIn order to improve the function of holding the magnets, an interference of the holding member may be increased so as to increase an amount of extension in the circumferential direction. However, if in particular, the holding member is made of FRP formed from a sheet-like material wound around a cylindrical jig, a terminal end portion of the sheet material or a portion nearby may not be able to extend enough that these portions may get loose or come off from the holding member, impairing the reliability of the rotor. In addition, as a result of increasing the interference of the holding member, great restorative force acts on the holding member, a fiber on an outer circumferential side may get into a gap between fibers on an inner circumferential side, and therefore the intended holding strength may be lost.\nAccording to a known art, the magnets and the protective pipe are held by the internal pressure of a filing material as disclosed in JP-A-2013-169103. In this case, a filling material is thermoplastic resin injected by injection molding. The internal pressure generated by the resin is small enough relative to tensile strength of the resin that the resin is not deformed. In other words, this known art can be only applied to the case where it is only necessary to provide fixing force which can be borne by plastic. If the rotor having such a structure is operated at high speed, the following problems arise. (1) The protective tube is expanded due to centrifugal force corresponding to the own weight of the protective tube, but the filling material cannot follow the expansion of the protective tube after reaching a certain limit. As a result, the fixing force acting on the protective tube is lost, and the protective tube starts to run idle. The protective tube scrapes the filling material with its inner wall, and is eventually detached from the rotor with noise. (2) As compared to the centrifugal force corresponding to the weight of the magnets having comparable density to iron, the fixing force by the filling material of resin is very weak, and may not be enough to fix the magnets against the centrifugal force. Therefore, when the rotor reaches a certain rotational speed, the magnets start to move to a radially outside and in a rotational direction, running idle and destroying the filling material. Accordingly, the structure disclosed in JP-A-2013-169103 is not appropriate to increase rotational speed.\nTherefore, there is a need for a rotor which can increase efficiency and reliability of an electric motor."}
{"text": "1. Field of the Invention\nThe present invention relates to an electron gun with band-shaped beams for producing an electron beam with a band-shaped cross-section. It is equipped with a linear cathode and is particularly suitable for low-energy electron beam systems. The electron beam exits into the atmosphere from a beam window, also called Lenard window. Preferred applications are the curing and cross-linkage of thin polymer coatings with a coating thickness of up to several 10 .mu.m on band- as well as slab-shaped coating substrates.\n2. Discussion of Background Information\nElectron beam devices are known, wherein short, wire-shaped cathodes that are stressed in the running direction of the product are arranged in its electron gun. These cathodes are arranged equidistant and parallel to one another at certain distances and individually suspended under tension by springs. Their number is determined by the respective beam widths and therefore the window length. The cathode length is derived from the window width. Each of these cathodes is concentrically surrounded by a cylindrical grid. In addition another broad grid which stretches over all the grid-cathode arrangements is generally provided. The grids have been assigned different potentials and serve to homogenize the current density laterally to the cathode voltage direction (U.S. Pat. No. 3,863,163).\nSince all of these cathodes are wired electrically in parallel, the disadvantage is that very high heating currents can arise--especially with large exposure widths and a correspondingly large number of parallel cathodes. Such cathode systems are used in particular for electron beam systems with acceleration voltages in the range of &gt;180 kV. The width of the beam window is thereby in the range of 200 mm. Another disadvantage of such electron beam systems is that if an outage of single cathodes occurs, the homogeneity of the current density in the window length is significantly disturbed. The change of a cathode is very costly and time-consuming due to the construction.\nElectron beam generators with a band-shaped beam cross-section whose beam window corresponds in length to the product width (DE 44 32 983 C1) are known. They have one or more parallel cathodes arranged over the window length and therefore the exposure width. A spring prestressing of the cathodes is necessary in order to limit the cathode sag. This cathode arrangement is surrounded by a generally tube-shaped control electrode, slitted in the beam direction, which serves to form the beam and if necessary to control the beam current.\nAll the devices have the disadvantage that with increasing window length, and thus also cathode length, a voltage strength is necessary which exceeds the yielding point of the cathode material at operating temperature. The result thereof is that with increased operating time the cathode length increases and the planned constructional oscillation path is reached. A cathode tear or an inadmissibly cathode sag leads to a critical limitation of the cathode's lifespan. The stretching of the cathode material thus limits the applicability of such cathode Systems based on cathode lengths, and therefore beam widths, of about one meter.\nIn addition, these cathode systems tend to undergo mechanical oscillations. Such oscillations can easily be instigated by external influences. The result of such cathode oscillations can be uncontrollable effects upon the beam formation and thus the homogeneity of the electron beam. The problems resulting from the oscillations also increase with the length of the voltage cathode system.\nTo limit the disadvantages connected with the voltage cathode system, one or more intermittent suspension devices of the cathode have been suggested. These solutions are of only limited use, since the heating of the cathodes is connected with a considerable change in length and the suspension must tolerate this change in length. To this is added an emission intrusion of the cathode at the place of suspension due to heat dissipation. A further disadvantage is that the suspension device itself has a local effect on the potential distribution near the cathodes, which can also have a negative impact on the homogeneity of the beam current density distribution."}
{"text": "The protection of crops from weeds and other vegetation which inhibit crop growth is a constantly recurring problem in agriculture. To help combat this problem, researchers in the field of synthetic chemistry have produced an extensive variety of chemicals and chemical formulations effective in the control of such unwanted growth. Chemical herbicides of many types have been disclosed in the literature and a large number are in commercial use.\nIn some cases, active herbicides have been shown to be more effective in combination than when applied individually. The result is often termed \"synergism\" since the combination demonstrates a potency or activity level exceeding that which it would be expected to have, based on a knowledge of the individual potencies of the components. The present invention resides in the discovery that certain substituted ureas and certain pyrrolidones, already known individually for their herbicidal potency, display this synergistic effect when applied in combination."}
{"text": "1. Field\nThis disclosure relates to wireless energy transfer, also referred to as wireless power transmission.\n2. Description of the Related Art\nEnergy or power may be transferred wirelessly using a variety of known radiative, or far-field, and non-radiative, or near-field, techniques. For example, radiative wireless information transfer using low-directionality antennas, such as those used in radio and cellular communications systems and home computer networks, may be considered wireless energy transfer. However, this type of radiative transfer is very inefficient because only a tiny portion of the supplied or radiated power, namely, that portion in the direction of, and overlapping with, the receiver is picked up. The vast majority of the power is radiated away in all the other directions and lost in free space. Such inefficient power transfer may be acceptable for data transmission, but is not practical for transferring useful amounts of electrical energy for the purpose of doing work, such as for powering or charging electrical devices. One way to improve the transfer efficiency of some radiative energy transfer schemes is to use directional antennas to confine and preferentially direct the radiated energy towards a receiver. However, these directed radiation schemes may require an uninterruptible line-of-sight and potentially complicated tracking and steering mechanisms in the case of mobile transmitters and/or receivers. In addition, such schemes may pose hazards to objects or people that cross or intersect the beam when modest to high amounts of power are being transmitted. A known non-radiative, or near-field, wireless energy transfer scheme, often referred to as either induction or traditional induction, does not (intentionally) radiate power, but uses an oscillating current passing through a primary coil, to generate an oscillating magnetic near-field that induces currents in a near-by receiving or secondary coil. Traditional induction schemes have demonstrated the transmission of modest to large amounts of power, however only over very short distances, and with very small offset tolerances between the primary power supply unit and the secondary receiver unit. Electric transformers and proximity chargers are examples of devices that utilize this known short range, near-field energy transfer scheme.\nTherefore a need exists for a wireless power transfer scheme that is capable of transferring useful amounts of electrical power over mid-range distances or alignment offsets. Such a wireless power transfer scheme should enable useful energy transfer over greater distances and alignment offsets than those realized with traditional induction schemes, but without the limitations and risks inherent in radiative transmission schemes."}
{"text": "1. Field of the Invention\nThis invention relates to semiconductor fabrication and more particularly to an improved process of globally planarizing a field dielectric by placing minimum mesas (hereinafter \"silicon risers\") in large area isolation regions.\n2. Description of the Relevant Art\nThe fabrication of an integrated circuit involves placing numerous devices in a single semiconductor substrate. Select devices are interconnected by a conductor which extends over a dielectric which separates or \"isolates\" those devices. Implementing an electrical path across a monolithic integrated circuit thereby involves selectively connecting isolated devices. When fabricating integrated circuits it must therefore be possible to isolate devices built into the substrate from one another. From this perspective, isolation technology is one of the critical aspects of fabricating a functional integrated circuit.\nA popular isolation technology used for an MOS integrated circuit involves the process of locally oxidizing silicon. Local oxidation of silicon, or LOCOS process involves oxidizing field regions between devices. The oxide grown in field regions are termed field oxide, wherein field oxide is grown during the initial stages of integrated circuit fabrication, before source and drain implants are placed in device areas or active areas. By growing a thick field oxide in isolation (or field) regions pre-implanted with a channel-stop dopant, LOCOS processing serves to prevent the establishment of parasitic channels in the field regions.\nWhile LOCOS has remained a popular isolation technology, there are several problems inherent with LOCOS. First, a growing field oxide extends laterally as a bird's-beak structure. In many instances, the bird's-beak structure can unacceptably encroach into the device active area. Second, the pre-implanted channel-stop dopant oftentimes redistributes during the high temperatures associated with field oxide growth. Redistribution of channel-stop dopant primarily affects the active area periphery causing problems known as narrow-width effects. Third, the thickness of field oxide causes large elevational disparities across the semiconductor topography between field and active regions. Topological disparities cause planarity problems which become severe as circuit critical dimensions shrink. Lastly, thermal oxide growth is significantly thinner in small field (i.e., field areas of small lateral dimension) regions relative to large field regions. In small field regions, a phenomenon known as field-oxide-thinning effect therefore occurs. Field-oxide-thinning produces problems with respect to field threshold voltages, interconnect-to-substrate capacitance, and field-edge leakage in small field regions between closely spaced active areas.\nMany of the problems associated with LOCOS technology are alleviated by an isolation technique known as the \"shallow trench process\". Despite advances made to decrease bird's-beak, channel-stop encroachment and non-planarity, it appears that LOCOS technology is still inadequate for deep submicron MOS technologies. The shallow trench process is better suited for isolating densely spaced active devices having field regions less than one micron in lateral dimension.\nThe trench process involves the steps of etching a silicon substrate surface to a relatively shallow depth, e.g., between 0.2 to 0.5 microns, and then refilling the shallow trench with a deposited dielectric. Some trench processes include an interim step of growing oxide on trench walls prior to the trench being filled with a deposited dielectric. After the trench is filled, it is then planarized to complete the isolation structure.\nThe trench process eliminates bird's-beak and channel-stop dopant redistribution problems. In addition, the isolation structure is fully recessed, offering at least a potential for a planar surface. Still further, field-oxide thinning in narrow isolation spaces does not occur and the threshold voltage is constant as a function of channel width.\nWhile the trench isolation process has many advantages over LOCOS, it cannot in all instances achieve complete global planarization across the entire semiconductor topography. The upper surface of fill dielectric in large isolation areas are generally at lower elevational levels than the upper surface fill dielectric in small isolation areas. The fill dielectric readily deposits in small area trenches such that the elevation of the fill dielectric in a small area is greater than the elevation within a large area trench. Accordingly, subsequent processing is needed to bring the large area fill topography to the same elevational level as small area fill topography.\nMost researchers have focused upon fairly complex processes for planarizing large and small area fill dielectrics. Those processes generally involve rework of the fill dielectric. A popular rework technique involves depositing a sacrificial layer across the fill dielectric topography, and then removing the sacrificial layer at the same etch rate as the underlying dielectric. Generally, the sacrificial layer is deposited as a low-viscosity liquid. Baking the liquid, or exposing it to ultraviolet light, causes the liquid to convert to solid form, as a solid-gel reaction. Popular sacrificial materials include photoresist, polymide or spin-on glass (SOG). The sacrificial layers generally etch back in a plasma until the topmost regions of the fill dielectric are exposed. The etch chemistry is then modified so that the sacrificial layer material and the underlying fill dielectric are etched at approximately the same rate. The etch is continued until all of the sacrificial layer has been etched away, leaving a somewhat planar dielectric upper surface.\nThe sacrificial etchback technique is generally valid only for the planarization of topographies in which features are less than 10 .mu.m (microns) apart. For large regions between trenches, the step height will not be reduced, since the photoresist thickness on top of such features will be the same as the thickness over the adjacent trench.\nIn an effort to overcome the shortcomings of the etchback process, a planarization block mask may be used. In this procedure, a liquid material is applied and developed as a planarization film followed by a block mask used to expose and develop this film. The block mask protects topography in wide, low regions from a subsequent etch plasma. The upper surface of high regions can then be removed to an elevational level commensurate with the protected low regions. The planarization block mask involves an additional lithography step and a mask which is produced by selectively reversing the mask used in producing the underlying topography. Mask reversal may involve errors due to changes in the resist thickness caused by the underlying pattern, misalignment, etc.\nIn an effort to eliminate the complex deposition, patterning and etch processes involved with sacrificial deposition, etch back, and the block masking, many manufacturers have directed their attention to chemical-mechanical polishing (CMP). Application of a chemical slurry and an abrasive polishing pad across the entire semiconductor topography, allows planarization of that topography commensurate with the planarity of the pad surface. Unfortunately, however, when force is applied to a pad, the pad will oftentimes conform to the unevenness of that topography. Thus, while high elevation areas, or peaks, receive substantial polishing, low elevational areas (or valleys) are also abraded.\nA better understanding of the problems inherent with CMP are illustrates in reference to FIGS. 1 and 2. FIG. 1 depicts a partial cross-section of a semiconductor topography 10. Topography 10 includes a substrate 12 which has been fashioned with a small area trench 14 and a large area trench 16 according to the trench isolation process. Trenches 14 and 16 define field areas between active regions, wherein active regions are defined as silicon mesas 18 which extend from an elevation equivalent to the trench base. Deposited upon and between mesas 18 is a fill dielectric 20.\nFIG. 2 illustrates a processing step subsequent to that shown in FIG. 1. Specifically, the upper surface of fill dielectric material 20 receives CMP. The polishing pad inherently flexes or conforms under pressure to the upper surface of dielectric 20, causing the polishing pad to attack and remove dielectric 20 upper surface in large area trench 16, albeit to a lesser degree than the removal of dielectric 20 in small area trench 14. If dielectric in large area trench 16 is not sufficiently thick to withstand the attack, the dielectric upper surface will be removed below the desired planar elevation. CMP causes extensive erosion of isolated silicon mesas (i.e., the silicon surface). Thus, a slight recess of dielectric 20 upper surface occurs at the conclusion of CMP. That recess is shown as reference numeral 22.\nA need therefore exists in producing a process which can planarize fill dielectric within large area trenches to an elevational level equal to that of fill dielectric within small area trenches. Large area trenches are those defined as having a width or length greater than 5 microns. The desired process must be simple to use. Accordingly, the desired process must not entail additional deposition, etch, and/or lithography steps attributable to the aforementioned etchback and/or block masking processes. The desired process would preferably involve a fill dielectric conditioned with an upper topography which can be blanket planarized."}
{"text": "1. Field of the Invention\nThis invention relates to plastic articles and more particularly relates to medical articles having improved gas and liquid vapor barriers.\n2. Description of the Related Art\nPolypropylene (PP) has long been used in molding and extruding operations for articles such as plastic medical containers and films for the food packaging industry. Polyethylene terephthalate (PET) has more recently been used in molding and extruding operations for these articles. However, PP and PET are somewhat permeable to nitrogen, oxygen, and other gases and vapors. As a result, PP and PET containers are inherently subject to transmission of gases. As the medical industry begins to place increased emphasis on the use of plastic medical products, these permeability problems have become more acute.\nIn particular, evacuated blood collection tubes must meet certain performance standards. Such performance standards generally include the ability to maintain greater than about 90% original draw volume over a one year period, and gas permeability clearly is detrimental to this need. Moreover, materials must also be capable of being sterilized by radiation, and substantially avoid interfering with tests and analysis. Thus, materials for such containers not only must resist gas and liquid vapor permeability problems, but they must also meet several other requirements.\nVarious techniques have therefore been devised in an attempt to reduce gas and vapor permeability of containers fabricated from PP, PET and other resins. Such techniques include addition of inorganic fillers, coating the containers with resins having barrier properties, plasma chemical vapor deposition coating of inorganic materials, and blending, laminating or co-extruding the resins with barrier resins.\nWhile such efforts have offered some improvement, the need to consistently meet high performance standards demands further improvement."}
{"text": "1. Field of the Invention\nThe present invention relates to lighting circuits, more particularly, the present invention relates to LED (Light Emitting Diode) lighting circuits.\n2. Description of the Related Art\nLEDs (Light Emitting Diodes) are recently replacing traditional incandescent and fluorescent illuminating devices as main lighting sources in many applications such as automobiles and home appliances because of their long lifespan, high optic efficiency, and low profile, etc.\nTraditional arts of LED dimming control are generally achieved by adjusting the forward current flowing through the LED. Taking a white-light LED for instance, its color temperature will become lower when the forward current flowing through it becomes smaller than its regular forward current. The aforementioned color temperature variance is not desired by the industry. Therefore, there is a need to provide a LED dimming control with stable color temperature performance."}
{"text": "This invention relates to internal combustion engines, and more particularly to exhaust valve constructions for use in such engines.\nInternal combustion engines can be operated on a wide variety of fuels and while most are run on either diesel fuel or gasoline, many are designed to run on bottled gas, such as liquid propane. Lift trucks, for example, are frequently fueled with bottled gas.\nHeretofore, special provisions have had to be made to adapt moving parts of internal combustion, reciprocating engines for use with liquid gas fuels, as liquid propane. Consequently, assembly problems have been accentuated due to the need for two sets of parts, one for gasoline and one for bottled gas.\nThe problem is being accentuated through the increased emphasis on low lead or lead-free gasolines which pose the need for still further modifications to component configurations to optimize efficiency."}
{"text": "The inventive concept relates generally to electronic memory technologies. More particularly, the inventive concept relates to a volatile memory device having reconfigurable refresh timing.\nA volatile memory device loses stored data in the absence of applied power. One potential cause of this data loss is charge leakage from storage elements. For example, in a typical dynamic random access memory (DRAM), each memory cell comprises an access switch and a storage capacitor. Charges may leak from the storage capacitor due to leakage current in the access switch; as an example, a metal oxide transistor (MOS) forming the access switch may experience leakage current through a PN junction. The loss of charges on the storage capacitor may change the stored data, which can produce errors or malfunctions during system operation.\nTo prevent data loss, DRAM devices perform periodic refresh operations on the stored data. A refresh operation of a DRAM cell typically involves detecting stored data from the cell and applying charges to the cell according to the detected data.\nRefresh operations of a DRAM device can be performed in various alternative modes, such as auto refresh mode or a self refresh mode, for example. In the self refresh mode, the DRAM device typically performs refresh operations on a group of sequentially addressed memory cells by sequentially varying an internal address in response to a refresh instruction signal.\nEach memory cell is typically refreshed periodically according to a predetermined refresh cycle time (tREF). For example, if a first refresh operation occurs at a time t=0, a second refresh operation may occur at a time t=tREF, a third refresh operation may occur at a time t=2*tREF, and so on.\nIn general, the refresh cycle time is determined by a data retention time, which can vary between different devices or memory cells. One source of this variation includes differences in process, voltage, and temperature (PVT) conditions of the different devices or memory cells. Another source of this variation includes heat applied during packaging or set assembling or other environmental factors, which is referred to as a variable retention time (VRT). In general, the refresh characteristics tend to deteriorate during packaging or set assembling, so a refresh cycle time margin in a test stage is typically set to a small value to provide a margin for change.\nDue to subtleties of the above sources of variation, it may be difficult to predict with precision the data retention time of all memory cells in a DRAM device. Consequently, it may be difficult to set the refresh cycle time of a DRAM device in a way that avoids unnecessary refresh operations without losing stored data.\nIn view of the foregoing, there is a general need for improved techniques to address deterioration of refresh characteristics in volatile memory devices such as DRAMs. In particular, there is a need to address deterioration that may occur through various stages of device fabrication, such as a post package repair and a post set assembly repair."}
{"text": "Fast progress in information technologies triggers the need for new materials for high-density optical memory storage (see for example Albota, M. et al., Science 1998 281, 1653; C. E. Olson, M. J. R. Previte, J. T. Fourkas. Nature Materials 1, 2002, 225; S, Kawata, Y. Kawata, Chem. Rev. 2000, 100, 1777; Photoreactive Materials for Ultrahigh-Density Optical Memory, Ed. M. Irie, Elsevier, Amsterdam, 1994; A. Renn,; U. P. Wild, A. Rebane. Chem. Rev. 2000, 100, 1741; D. S. Tyson, C. A. Bignozzi, F. N. Castellano, J. Am. Chem. Soc. 2002, 124, 4562; and J. R. Sheats, P. F. Barbara, Acc. Chem. Res. 1999, 32, 191).\nAn increase in storage capacity can be achieved by shifting from two-dimensional to three-dimensional (3D) optical data storage: holographic recording with photorefractive media (Photorefarctive materials and Their Applications I, edited by P. Gunter and J.-P. Huignard (Springer, Berlin, 1988)), spectral hole burning (Persistent Spectral Hole Burning: Science and Applications, edited by W. E. Moerner (Springer, Berlin, 1987)) and photon echo (R. Kachru and M. K. Kim, Opt. Lett 28, 2186 (1989)).\nAn increase in storage density has been further accomplished by the use of two-photon processes introduced by Rentzepis (D. A. Parthenopoulos and P. M. Rentzepis, Science 245, 843 (1989); A. S. Dvornikov and P. M. Rentzepis, Opt. Commun. 119, 341 (1995)) and Webb, (W. Denk, J. H. Strickler, W. W. Webb, Science 248, 73 (1990)). On the materials side, new modes in data storage have been examined by employing polymer photonic crystals (B. Siwick, O. Kalinina, E. Kumacheva, R. J. D. Miller, J. Noolandi, J. Appl. Phys. 90 5328 (2001); N. I. Koroteev, S. A. Krikunov, S. A. Magnitskii, D. V. Malakhov, V. V. Shubin, Jpn. J. Appl. Phys. Part 1 37, 2279-2280 (1998), and D. Kraemer, B. Siwick, Miller R. J. D. J. Chem. Phys. 285, 73 (2002). These crystals can be prepared via “bottom-to-top” approach from the core-shell latex particles with fluorescent cores and optically inert shells (Kalinina, O., Kumacheva, E. Macromolecules 32, 4122-4129 (1999); Kumacheva, E., Kalinina, O., Lilge, L. Adv. Mater. 11, 231-234 (1999)). Encryption or recoding of information in 3D polymer photonic crystals was achieved by local one-photon or two-photon photobleaching of the fluorescent dye localized in periodic domains (bits). The existence of the “dead” space between the bits, acting as a barrier to cross-talk, led to a two-fold increase of signal-to-noise ratio compared to the material with a uniform dye distribution.\nMoreover, information retrieval with an order of magnitude beyond the Rayleigh limit was recently predicted due to photobleaching of domains with well-defined Fourier components (Kraemer D. Kraemer, B. Siwick, Miller R. J. D. J. Chem. Phys. 285, 73 (2002)).\nThe use of bicolored (or multicolored) multiphase periodic medium for bit-like optical data storage has several advantages. First, in the bit-like binary memory storage, reading is achieved by distinguishing between “1” s (photobleached domains) and “0” (fluorescent domains) (B. Siwick, O. Kalinina, E. Kumacheva, R. J. D. Miller, J. Noolandi, J. Appl. Phys. 90 5328 (2001). Thus reading of optically-inert (non-fluorescent) matrix is similar to reading of “1” s; this uncertainty may induce errors in retrieving information. The incorporation of the second dye in the matrix of the photonic crystal resolves this problem. More important, the use of several dyes with distinct excitation and emission peaks increases the number of modes in data storage from 2 (binary encryption mode, one-dye system) to 2n (binary code, multi-dye system), where n is the number of dyes. The incorporation of different dyes in different phases of the material reduces energy transfer between them. Data recording achieved by photo-induced changes in different dyes accompanied by information retrieval at a particular well-defined wavelength tailors additional security features to the recorded information.\nPatent Publication WO 03/044276 discloses a procedure for preparation of a paper containing security elements (photoluminescent objects such as fiber, thread, rod, tape, film, windows and combination of thereof) by their incorporation into paper structure. Authenticity of the paper is then determined using polarized light.\nJP 2002317134 discloses a fluorescent ink composition comprising a non-fluorescent dye absorbing in a visible range and fluorescent dye with no absorption in the visible range and fluorescence in the visible range. Authenticity of the document is determined by irradiating it with a light source capable of exciting fluorescent dye and observing an emission.\nJP 2002088688 discloses a composition for laser printing with toner particles mixed with or coated by fluorescent dye or pigment absorbing in UV range. Authenticity of the document is checked by exciting the UV dye and observing fluorescence.\nWO 01/09435 discloses a security paper comprising at least two of the following security features: a water insoluble dye which is incorporated in the paper and bleeds when placed in contact with organic solvent; a dye that becomes fluorescent in daylight when brought in contact with alkaline substance; a dye which becomes fluorescent when exposed to UV light.\nNone of above-mentioned patents teaches or suggests information encryption using multiple imaging and selective photobleaching. All features are used to prove authenticity of the document or paper whereas the material produced in accordance with the present invention can be used to identify a document holder and demonstrate the authenticity of the document.\nUS2003/0183695 discloses a method for making a secure ID card with multiple images. The multicolored images are printed on the information-bearing layer in such a way as to provide multiple images when printed information is viewed at different predetermined angles through the lenticular lens. This method does not use selective photobleaching of fluorescent dyes. There is a high probability that the images can be viewed all together or one by one when different angles are used.\nPatent publication WO 03101755 discloses an article for authentication (e.g. security document or sheet made from a transparent polymeric film which comprises a first luminescent material in a first region and a second luminescent material in a second region. The second region contains a transparent window and the rest of the material is optionally opacified. The two materials luminesce at different visible wavelengths when excited (e.g. the coating fluoresces red/orange when excited with UV-C at 254 nm whereas the window fluoresces blue when excited with UV-A at 365 nm). However when the article is folded the first region is viewed through the transparent window and when excited the combined luminescence of both materials is seen as a colour change which may be used to authenticate the article. This method provides only one authentication feature: the color of fluorescence which is not easy to identify without expensive equipment.\nIt would be particularly advantageous to provide an optical data storage device and a method of writing or storing data in the optical storage device using selective single-photon or two-photon photobleaching."}
{"text": "Many computing systems (for example, portable or mobile systems) need to maintain a low power consumption. For example, net-zero energy computing devices, wearable devices, or internet of things (IoT) devices, among others, require low power consumption. Net-zero energy computing devices are experiencing a paradigm shift toward extremely low power consumption. In particular, reduction of standby power can be an important goal since some devices may spend most of their time in a standby mode. In order to maintain a low overhead, it is key for some devices to transition very quickly between low power standby modes and active modes. However, some circuits (for example, some analog circuits) that are designed for very low power consumption may exhibit low speeds. Reference voltage generators that define reference voltages for multiple supply rails of an integrated circuit (IC) such as an SoC (system on a chip) can exhibit low speeds. Reference voltages may need to be maintained with very low power consumption while in a standby state, but may also need to quickly change (for example, when an SoC transitions to an active state and requires higher supply voltages).\nIn some cases, the same numbers are used throughout the disclosure and the figures to reference like components and features. In some cases, numbers in the 100 series refer to features originally found in FIG. 1; numbers in the 200 series refer to features originally found in FIG. 2; and so on."}
{"text": "In the field of all solid state batteries, attention has been paid to the interface of an active material and a solid electrolyte material, and there have been attempts to promote an enhancement in the performance of batteries. For example, Patent Literature 1 discloses an all solid lithium battery including a lithium ion conductive solid electrolyte having a sulfide as a main component and a cathode active material having a surface coated with a lithium ion conductive oxide. Further, as an example of the lithium ion conductive oxide, LiNbO3 is disclosed. Furthermore, as an example of a method for forming a lithium ion conductive oxide, a method of using a solution containing a metal alkoxide is disclosed.\nFurthermore, Patent Literature 2 discloses a lithium-transition metal oxide powder having a surface coated with a coating layer containing lithium niobate and the lithium-transition metal oxide powder having a carbon content of 0.03% by mass or less. Furthermore, as an example of a method for forming lithium niobate, a method of using a solution containing a peroxo complex of niobium acid is disclosed."}
{"text": "Electrical transformers such as of a type used in commercial, power distribution and industrial applications are typically constructed with the transformer coils located in a housing that is filled with liquid to achieve desired electrical insulation and thermal cooling characteristics. Liquids used for this purpose include various oil formulations and hazardous chemicals. In many instances, leaking of these fluids from transformers can cause environmental damage, violate certain civil codes, and result in the levying of significant fines. Thus, there is an ever present need to prevent leaking from transformers into the outside environment.\nTransformers can start leaking for a number of common reasons, generally related to either damage to the transformer housing or time in service, and as such, development of leaks are often not preventable. By the time maintenance personnel reaches a leaking transformer, the leaking of fluid will have typically stopped. The significant problem for the maintenance personnel is to remove the leaking transformer and return it to a repair facility without leaking additional fluid. Thus, a cause of transformers needlessly leaking fluid into the outside environment is during the course of handling a leaking transformer, including removal of the installed transformer and transport of the transformer to a repair facility.\nPrior transformer containment bags provide no capability of lifting or carrying or moving the transformer. The transformer cannot be lifted or picked up and moved with the containment bag. In order to move a transformer, the lift device must connect to the transformer. Thus, even if the containment bag is secured to the transformer, portions of the transformer are not contained within the bag during movement of the transformer by the lift device, which presents the danger of further leaking from the transformer to the outside environment when transformer is, for example, lowered onto a pallet or truck for transport to a storage or repair facility. This danger is enhanced where the leak is located above the lift locations on the transformer, such as when bushing which are located above the lift locations of pole mount transformers leak due to time in service.\nDuring transport to a storage or repair facility, such as in the back of a repair truck, the transformers can experience significant bouncing and “knocking-around” during which prior transformer containment bags provide no protection.\nThe present invention addresses the above-identified and other known drawbacks and disadvantages of prior transformers containment bags."}
{"text": "Data driven applications include tools developed to access information within a database using a predefined data model design specific to the particular application (and perhaps also specific to customer requirement(s)). As these applications are revised, upgraded, and/or superseded the data model design for the application can change to incorporate these revisions, upgrades, and changes. When changes to the application software introduce deviations in the data model design, the conventional database can undergo an update of its database structures to remain compatible (e.g., accessible, readable, writeable, etc.) with the application. Database upgrade and/or modification can consume a great amount of time and system resources proportional to the magnitude of the data set(s) within the database. Such upgrade/modification can run the risk of data loss. This loss can be permanent as roll back of the database structure can often not be possible."}
{"text": "Traditionally, leases of well services involve a leassor creating a lease agreement with a leassee. In creating the lease relationship and lease agreement, the prospective leassee will generally require information regarding the well services from the leassor. At times, leassors may lack the time or expertise to fully monitor the well services that are being provided, with precision and regularity, thereby creating an inability for them to thoroughly inform the leassees of the information that they require.\nSome leassors turn to management companies. However, these companies can be expensive, and there are no assurances that payments to the leassor will be made when the payments are due. Additionally, these property management companies can lack the resources and expertise to provide the requisite information that a potential leassee requires before leasing well services.\nAccordingly, there exists a need for a system for acquiring rights to lease well services for a specified period of time or a particular activity.\nA potential leassee and leassor will normally require certain information before actually leasing well services, such as the availability of the well services, any technical specifications regarding the well services, the number of bid requests a potential leassee will have, contract terms of the lease, the different types of rates for services offered by the leassor, and the different types of responsibilities the leassee and leassor will incur. The enormous amount of information that needs to transfer between the potential leassee and leassor increases the probability of delaying, misplacing, or forgetting a piece of information that would be vital to either party.\nThus, a need exists to provide a system for a plurality of potential leassees to obtain information regarding a plurality of well services. A need exists to provide a system for a plurality of leassors to obtain information regarding the potential leassees that wish to lease well services. Further, a need exists for a database server system to interface the needs of the potential leassee and the leassor in a manageable fashion.\nThe present embodiments meet these needs.\nThe present embodiments are detailed below with reference to the listed Figures."}
{"text": "In a conventional reciprocating piston, internal combustion engine, torque is near zero at low engine output speed and peaks somewhere in the middle of the engine operating range and then decreases as output speed further increases. On the other hand, a gas turbine engine has very high torque at low engine speeds.\nTo provide a friction start clutch in a transmission for a high torque engine, such as a gas turbine engine, one would need a very wide range of torque capacity, and thus, a very high capacity clutch. A single large clutch piston would require a clutch piston apply area having sensitivity at a high level. If 100 psi were necessary to apply the clutch, a 1.0 psi variation in pressure would result in a large torque variation which would be objectionable at light throttle starts. Such an adjustment to provide adequate clutch apply pressure would not provide a controllable sensitivity.\nThe present invention provides a novel clutch arrangement and control strategy for transmission starting in a high torque engine."}
{"text": "The present invention relates to compositions for polymerizing one or more monomers or mixtures of monomers such as ethylene and one or more comonomers, to form an interpolymer product having unique physical properties, to a process for preparing such interpolymers, and to the resulting polymer products. In another aspect, the invention relates to methods of using these polymers in applications requiring unique combinations of physical properties. In still another aspect, the invention relates to the articles prepared from these polymers. The inventive polymers comprise two or more differing regions or segments (blocks) causing the polymer to possess unique physical properties. These multi-block copolymers and polymeric blends comprising the same are usefully employed in the preparation of solid articles such as moldings, films, sheets, and foamed objects by molding, extruding, or other processes, and are useful as components or ingredients in adhesives, laminates, polymeric blends, and other end uses. The resulting products are used in the manufacture of components for automobiles, such as profiles, bumpers and trim parts; packaging materials; electric cable insulation, and other applications.\nIt has long been known that polymers containing a block-type structure often have superior properties compared to random copolymers and blends. For example, triblock copolymers of styrene and butadiene (SBS) and hydrogenated versions of the same (SEBS) have an excellent combination of heat resistance and elasticity. Other block copolymers are also known in the art. Generally, block copolymers known as thermoplastic elastomers (TPE) have desirable properties due to the presence of “soft” or elastomeric block segments connecting “hard” either crystallizable or glassy blocks in the same polymer. At temperatures up to the melt temperature or glass transition temperature of the hard segments, the polymers demonstrate elastomeric character. At higher temperatures, the polymers become flowable, exhibiting thermoplastic behavior. Known methods of preparing block copolymers include anionic polymerization and controlled free radical polymerization. Unfortunately, these methods of preparing block copolymers require sequential monomer addition and batch processing and the types of monomers that can be usefully employed in such methods are relatively limited. For example, in the anionic polymerization of styrene and butadiene to form a SBS type block copolymer, each polymer chain requires a stoichiometric amount of initiator and the resulting polymers have extremely narrow molecular weight distribution, Mw/Mn, preferably from 1.0 to 1.3. Additionally, anionic and free-radical processes are relatively slow, resulting in poor process economics.\nIt would be desirable to produce block copolymers catalytically, that is, in a process wherein more than one polymer molecule is produced for each catalyst or initiator molecule. In addition, it would be highly desirable to produce block copolymers from olefin monomers such as ethylene, propylene, and higher alpha-olefins that are generally unsuited for use in anionic or free-radical polymerizations. In certain of these polymers, it is highly desirable that some or all of the polymer blocks comprise amorphous polymers such as a copolymer of ethylene and a comonomer, especially amorphous random copolymers comprising ethylene and an α-olefin having 3, and especially 4, or more carbon atoms. Finally, if would be highly desirable to be able to use a continuous process for production of block copolymers.\nPrevious researchers have stated that certain homogeneous coordination polymerization catalysts can be used to prepare polymers having a substantially “block-like” structure by suppressing chain-transfer during the polymerization, for example, by conducting the polymerization process in the absence of a chain transfer agent and at a sufficiently low temperature such that chain transfer by β-hydride elimination or other chain transfer processes is essentially eliminated. Under such conditions, the sequential addition of different monomers was said to result in formation of polymers having sequences or segments of different monomer content. Several examples of such catalyst compositions and processes are reviewed by Coates, Hustad, and Reinartz in Angew. Chem., Int. Ed., 41, 2236-2257 (2002) as well as US-A-2003/0114623.\nDisadvantageously, such processes require sequential monomer addition and result in the production of only one polymer chain per active catalyst center, which limits catalyst productivity. In addition, the requirement of relatively low process temperatures increases the process operating costs, making such processes unsuited for commercial implementation. Moreover, the catalyst cannot be optimized for formation of each respective polymer type, and therefore the entire process results in production of polymer blocks or segments of less than maximal efficiency and/or quality. For example, formation of a certain quantity of prematurely terminated polymer is generally unavoidable, resulting in the forming of blends having inferior polymer properties. Accordingly, under normal operating conditions, for sequentially prepared block copolymers having Mw/Mn of 1.5 or greater, the resulting distribution of block lengths is relatively inhomogeneous, not a most probable distribution. Finally, sequentially prepared block copolymers must be prepared in a batch process, limiting rates and increasing costs with respect to polymerization reactions carried out in a continuous process.\nFor these reasons, it would be highly desirable to provide a process for producing olefin copolymers in well defined blocks or segments in a process using coordination polymerization catalysts capable of operation at high catalytic efficiencies. In addition, it would be desirable to provide a process and resulting block or segmented copolymers wherein insertion of terminal blocks or sequencing of blocks within the polymer can be influenced by appropriate selection of process conditions. Finally, it would be desirable to provide a continuous process for producing multi-block copolymers.\nThe use of certain metal alkyl compounds and other compounds, such as hydrogen, as chain transfer agents to interrupt chain growth in olefin polymerizations is well known in the art. In addition, it is known to employ such compounds, especially aluminum alkyl compounds, as scavengers or as cocatalysts in olefin polymerizations. In Macromolecules, 33, 9192-9199 (2000) the use of certain aluminum trialkyl compounds as chain transfer agents in combination with certain paired zirconocene catalyst compositions resulted in polypropylene mixtures containing small quantities of polymer fractions containing both isotactic and atactic chain segments. In Liu and Rytter, Macromolecular Rapid Comm., 22, 952-956 (2001) and Bruaseth and Rytter, Macromolecules, 36, 3026-3034 (2003) mixtures of ethylene and 1-hexene were polymerized by a similar catalyst composition containing trimethylaluminum chain transfer agent. In the latter reference, the authors summarized the prior art studies in the following manner (some citations omitted):                “Mixing of two metallocenes with known polymerization behavior can be used to control polymer microstructure. Several studies have been performed of ethene polymerization by mixing two metallocenes. Common observations were that, by combining catalysts which separately give polyethene with different Mw, polyethene with broader and in some cases bimodal MWD can be obtained. [S]oares and Kim (J. Polym. Sci., Part A: Polym. Chem., 38, 1408-1432 (2000)) developed a criterion in order to test the MWD bimodality of polymers made by dual single-site catalysts, as exemplified by ethene/1-hexene copolymerization of the mixtures Et(Ind)2ZrCl2/Cp2HfCl2 and Et(Ind)2ZrCl2/CGC (constrained geometry catalyst) supported on silica. Heiland and Kaminsky (Makromol. Chem., 193, 601-610 (1992)) studied a mixture of Et-(Ind)2ZrCl2 and the hafnium analogue in copolymerization of ethene and 1-butene.        These studies do not contain any indication of interaction between the two different sites, for example, by readsorption of a terminated chain at the alternative site. Such reports have been issued, however, for polymerization of propene. Chien et al. (J. Polym. Sci., Part A: Polym. Chem., 37, 2439-2445 (1999), Makromol., 30, 3447-3458 (1997)) studied propene polymerization by homogeneous binary zirconocene catalysts. A blend of isotactic polypropylene (i-PP), atactic polypropylene (a-PP), and a stereoblock fraction (1-PP-b-a-PP) was obtained with a binary system comprising an isospecific and an aspecific precursor with a borate and TIBA as cocatalyst. By using a binary mixture of isospecific and syndiospecific zirconocenes, a blend of isotactic polypropylene (i-PP), syndiotactic polypropylene (s-PP), and a stereoblock fraction (1-PP-b-s-PP) was obtained. The mechanism for formation of the stereoblock fraction was proposed to involve the exchange of propagating chains between the two different catalytic sites. Przybyla and Fink (Acta Polym., 50, 77-83 (1999)) used two different types of metallocenes (isospecific and syndiospecific) supported on the same silica for propene polymerization. They reported that, with a certain type of silica support, chain transfer between the active species in the catalyst system occurred, and stereoblock PP was obtained. Lieber and Brintzinger (Macromol. 3, 9192-9199 (2000)) have proposed a more detailed explanation of how the transfer of a growing polymer chain from one type of metallocene to another occurs. They studied propene polymerization by catalyst mixtures of two different ansa-zirconocenes. The different catalysts were first studied individually with regard to their tendency toward alkyl-polymeryl exchange with the alkylaluminum activator and then pairwise with respect to their capability to produce polymers with a stereoblock structure. They reported that formation of stereoblock polymers by a mixture of zirconocene catalysts with different stereoselectivities is contingent upon an efficient polymeryl exchange between the Zr catalyst centers and the Al centers of the cocatalyst.”        \nBrusath and Rytter then disclosed their own observations using paired zirconocene catalysts to polymerize mixtures of ethylene/1-hexene and reported the effects of the influence of the dual site catalyst on polymerization activity, incorporation of comonomer, and polymer microstructure using methylalumoxane cocatalyst.\nAnalysis of the foregoing results indicate that Rytter and coworkers likely failed to utilize combinations of catalyst, cocatalyst, and third components that were capable of readsorption of the polymer chain from the chain transfer agent onto both of the active catalytic sites, i.e., two-way readsorption. While indicating that chain termination due to the presence of trimethylaluminum likely occurred with respect to polymer formed from the catalyst incorporating minimal comonomer, and thereafter that polymeryl exchange with the more open catalytic site followed by continued polymerization likely occurred, evidence of the reverse flow of polymer ligands appeared to be lacking in the reference. In fact, in a later communication, Rytter, et. al., Polymer, 45, 7853-7861 (2004), it was reported that no chain transfer between the catalyst sites actually took place in the earlier experiments. Similar polymerizations were reported in WO98/34970.\nIn U.S. Pat. Nos. 6,380,341 and 6,169,151, use of a “fluxional” metallocene catalyst, that is a metallocene capable of relatively facile conversion between two stereoisomeric forms having differing polymerization characteristics such as differing reactivity ratios was said to result in production of olefin copolymers having a “blocky” structure. Disadvantageously, the respective stereoisomers of such metallocenes generally fail to possess significant difference in polymer formation properties and are incapable of forming both highly crystalline and amorphous block copolymer segments, for example, from a given monomer mixture under fixed reaction conditions. Moreover, because the relative ratio of the two “fluxional” forms of the catalyst cannot be varied, there is no ability, using “fluxional” catalysts, to vary polymer block composition or the ratio of the respective blocks. Finally, prior art methods for olefin block copolymerization have been incapable of readily controlling the sequencing of the various polymer blocks, and in particular controlling the nature of the terminating block or segment of a multi-block copolymer. For certain applications, it is desirable to produce polymers having terminal blocks that are highly crystalline, that are functionalized or more readily functionalized, or that possess other distinguishing properties. For example, it is believed that polymers wherein the terminal segments or blocks are crystalline or glassy possess improved abrasion resistance and thermal properties such as tensile strength, elastic recovery and compression set. In addition, polymers wherein the blocks having amorphous properties are internal or primarily connected between crystalline or glassy blocks, have improved elastomeric properties, such as improved retractive force and recovery, particularly at elevated temperatures.\nIn JAGS, 2004, 126, 10701-10712, Gibson, et al discuss the effects of “catalyzed living polymerization” on molecular weight distribution. The authors define catalyzed living polymerization in this manner:\n“ . . . if chain transfer to aluminum constitutes the sole transfer mechanism and the exchange of the growing polymer chain between the transition metal and the aluminum centers is very fast and reversible, the polymer chains will appear to be growing on the aluminum centers. This can then reasonably be described as a catalyzed chain growth reaction on aluminum . . . . An attractive manifestation of this type of chain growth reaction is a Poisson distribution of product molecular weights, as opposed to the Schulz-Flory distribution that arises when β-H transfer accompanies propagation.”\nThe authors reported the results for the catalyzed living homopolymerization of ethylene using an iron containing catalyst in combination with ZnEt2, ZnMe2, or Zn(i-Pr)2. Homoleptic alkyls of aluminum, boron, tin, lithium, magnesium and lead did not induce catalyzed chain growth. Using GaMe3 as cocatalyst resulted in production of a polymer having a narrow molecular weight distribution. However, after analysis of time-dependent product distribution, the authors concluded this reaction was, “not a simple catalyzed chain growth reaction.” The reference fails to disclose the use of two or more catalysts in combination with a chain shuttling agent to make multi-block copolymers. Similar processes employing single catalysts have been described in U.S. Pat. Nos. 5,210,338, 5,276,220, and 6,444,867.\nEarlier workers have claimed to have formed block copolymers using a single Ziegler-Natta type catalyst in multiple reactors arranged in series, see for example U.S. Pat. Nos. 3,970,719 and 4,039,632. Additional Ziegler-Natta based processes and polymers are disclosed in U.S. Pat. Nos. 4,971,936; 5,089,573; 5,118,767; 5,118,768; 5,134,209; 5,229,477; 5,270,276; 5,270,410; 5,294,581; 5,543,458; 5,550,194; and 5,693,713, as well as in EP-A-470,171 and EP-A-500,530.\nDespite the advances by the foregoing researchers, there remains a need in the art for a polymerization process that is capable of preparing block like copolymers, especially multi-block copolymers, and most especially linear multi-block copolymers, in high yield and selectivity. Moreover, it would be desirable if there were provided an improved process for preparing multi-block copolymers, especially linear multi-block copolymers, of two or more olefin monomers such as ethylene and one or more comonomers, by the use of a shuttling agent. In addition it would be desirable to provide such an improved process that is capable of preparing multi-block copolymers, especially linear multi-block copolymers, having a relatively narrow molecular weight distribution. It would further be desirable to provide an improved process for preparing copolymers having more than two segments or blocks. Furthermore, it would be desirable to provide a process for identifying combinations of catalysts and chain shuttling agents capable of making such multi-block copolymers. Even further, it would be desirable to provide a process for independent control of the order of the various polymer blocks, especially a process for preparing olefin block copolymers containing terminal blocks having high crystallinity and/or functionality. Finally, it would be desirable to provide an improved process for preparing any of the foregoing desirable polymer products in a continuous process, without required sequential addition of monomers. Highly desirably, such process allows for independent control of the quantity and/or identity of the shuttling agent(s) and/or catalysts used."}
{"text": "1. Field of the Invention\nThe present invention relates to a sheet processing apparatus and image forming apparatus and more particularly to the sheet processing apparatus and image forming apparatus which perform a processing of binding a sheet bundle having a front cover and a back cover.\n2. Description of the Related Art\nIn recent years, diversification of information has been accelerated and accompanied by this, an opportunity for a person to output a content having a high confidentiality from a printer which is an example of the image forming apparatus individually and possess it has been increased. Such an output needs to be treated so that the content thereof cannot be seen easily by other people.\nTo solve this problem, there have been invented a plurality of methods for blocking a bookbound matter from being opened so that the content is hard to see by binding a fore edge, or an end opposite to a bound end of the bookbound matter temporarily (see Japanese Patent Application Laid-Open Nos. 2004-284750 and 9-165136).\nFIG. 18A illustrates an example of the method for blocking bookbound matter from being opened. Folded sheets S having perforations b at a bent portion are overlaid and the fore edge is bound with binding devices, for example, needles SP to make the bookbound matter double-leaved. In the meantime, this bookbound matter can be seen by separating across the perforations b.\nThe bookbound matter shown in FIG. 18B is double-leaved by binding the fore edge of sheet bundles SA each side-stitched, using the needles SP. This bookbound matter allows its content to be seen by removing the needles SP which bind the fore edges.\nHowever, in case of the conventional bookbound matter formed in this way, if the binding devices are removed to see the content thereof, it is difficult to return the state to the original bound condition. Particularly, in case of a bookbound matter having a large number of sheets, once the double-leaved sheets are released, it cannot be double-leaved again using a marketed stapler and if the fore edges are bound with the stapler again, a plurality of holes are made in the bookbound matter, thereby producing a bad appearance of the sheets.\nFurther, when any binding device is removed from the fore edge to release the double-leaved sheets, if no care is taken, the sheet might be torn out or bent.\nAccordingly, the present invention has been achieved in views of such a circumstance and the present invention intends to provide a sheet processing apparatus and image forming apparatus capable of producing double-leaved sheets and released that binding condition, without use of any binding device."}
{"text": "A number of applications exist wherein it is desirable to provide panels or boards having desirable acoustical and/or thermal properties. For example, it is desirable to provide molded panels for use as liners in heavy machinery for purposes of sound and/or thermal attenuation, such applications including molded liners for the interior of truck cabs and the like, liners for engine compartments and fire walls, and so forth. In many applications, it has been found that panels formed of fiberglass have yielded desirable results although other textile fabrics and fibers may be employed such as resinated cotton, synthetic fibers, reclaimed fibers and the like.\nIn applications wherein it is desirable to provide a predetermined level of sound attenuation, the desired attenuation may be obtained through the use of boards produced from fiberglass or other suitable fibrous media. The sound attenuation capability of such a fiberglass board is a function of the density and thickness of the board. One technique for producing boards having desirable thermal and/or acoustic properties is to mold the boards from basic blankets which are comprised of loosely gathered fiberglass filaments and a suitable binder. The fibers are preferably saturated with a resin or a binder such as a phenolic or thermoplastic or other thermosetting binders which are utilized, when cured, to hold the fibers together in order to assure the proper molded shape of the molded end product.\nIn order to obtain the desired density, it is conventional to utilize the aforementioned blankets of reasonably standardized density and to place the desired number of blankets within a mold in order to form a panel or board having the desired acoustic and thermal properties. For example, blankets of the type described may have a nominal uncompressed density of one pound per cubic foot and a nominal uncompressed thickness of one inch. Assuming it is desired to produce a board having a thickness of 0.5\" and a density of 8 lbs. per cubic foot, the desired board may be produced of four (4) of the aforementioned blankets which are inserted within a mold designed to compress the blankets to a total thickness of 0.50\". Insertion of four such blankets into the mold produces an end product having a density of 8 lbs. per cubic foot. The mold is normally heated and is maintained at a temperature level sufficient to cause the binder to set, thereby obtaining a board having the desired thickness and density. The blankets are retained in the mold for a period of time sufficient to set the binder and thereby provide a finished product having a substantially uniform density and thickness in order to achieve the desired acoustical and/or thermal attenuation properties.\nA number of applications exist wherein it is desired to provide acoustical panels and the like with predetermined shapes and with portions thereof having thicknesses which significantly deviate from the nominal thickness. For example, in applications wherein the molded members are designed for use as liners in the cabs of large industrial vehicles, the liners are typically designed to conform to the interior shape of the cab, such as, for example, the cab ceiling, and interior walls, hence the necessity for shaped molded members. In order to facilitate mounting of the liners to the interior surface of the cab, it is frequently desirable to provide openings for receiving fastening members and to provide the marginal portions of said openings with a thickness which is significantly reduced from the nominal thickness of the liners. In the example where a molded sheet formed of four standard blankets having densities of 1 lb. per cubic foot and uncompressed thicknesses of 1\", the portions of reduced thickness of the molded end product may be as small at 1/8\" making the density in such regions as high as 32 lbs. per cubic foot. The density and/or packing of the fibers in the regions of significantly reduced thickness is so great as to prevent the molds from compressing the original blankets down to the desired thickness, thereby requiring use of fewer blankets or blankets having lower nominal densities in order to allow the portions of the molded member of significantly reduced thickness to be molded to the desired thickness. This results in a reduced density for the major portion of the liner, which yields an undesirable compromise between the thickness and acoustical and/or thermal properties of the molded end product.\nIt is also required in some applications to provide certain surface finishes for the liner necessitating that the liner have a smooth high density finish capable of receiving and retaining the surface finish to be applied thereto. These requirements further increase the density requirements of the end product, which contributes to undesirable increases in material costs, process costs, labor costs, weight and brittleness.\nThe problems of panel brittleness result from the fact that the end product requires an increased density sufficient to provide a surface capable of being finished in the predetermined manner with the result that the end product is overly brittle and thus highly susceptible to breakage.\nSurface coatings or facings require a good high density surface to promote application and retention, thus requiring higher density panels which necessitate the use of a greater amount of fibrous material.\nThe criteria to be considered in forming molded parts of the type described above are, in the order of importance, the selection of the materials for achieving the desired thermal and/or acoustic properties; the structural requirements; and the shape and thickness of the molded part. The desired acoustic and/or thermal attenuation is a function of the thickness and density of the molded part. The ideal thickness of the sound absorbent material is one half the wave length of the sound to be absorbed. Assuming the sound to be absorbed has a frequency of 250 cycles per second, one half the wavelength of this frequency is 2.25 ft. which is clearly an impractical design from the viewpoint of both cost and size. Thus, conventional present-day designs are the result of a compromise which is struck between a group of the primary factors which contribute to the design of acoustic and/or thermal panels and the like.\nStructural requirements frequently necessitates changes in density to provide adequate supporting or mounting strength for portions of the panel.\nThe shape and thickness of the molded part impose limitations on maximum permissible density, necessitating reduction in overall density in order to be assured of achieving the desired molded shapes, the practical limitations of the raw material being such that the fibers of the blanket do not \"flow\" in the same manner as plastic material flows in order to freely and easily assume desired shapes and configurations.\nThe interplay of the three above-mentioned criteria result in a molded end product having an overall design in which density and thickness are compromised to achieve a satisfactory balance of the desirable operating and design characteristics.\nIt therefore becomes desirable to provide a novel method and apparatus for producing molded members and novel molded members having desirable acoustical and/or thermal properties and which more favorably conform to the important design criteria for such panels and which have improved characteristics as compared with panels produced through the use of conventional methods."}
{"text": "This invention pertains to a system for imparting notches into fiber optic cable trough components to facilitate their connection or attachment.\nIn the telecommunications industry there are numerous locations where a significant amount of fiber optic cable must be routed within a facility or from one facility to another. The routing within a facility may be from one piece of equipment to another or even from outside lines coming into a central office and to connectors where the lines are connected to equipment within the facility.\nThere are typically a large number of fibers and all the fibers must all be handled with care to avoid damage to the fiber optic cable, which would hinder its performance.\nIn a typical facility, fiber optic troughs are normally used to carry or route the fiber optic cables. In many facilities, the troughs or raceways are located overhead and over the locations of the fiber optic distribution frames, bays and equipment.\nA significant trough network or raceway configuration may be needed to contain and route the fiber optic cables. The installation time and expense can be substantial for trough systems and is further increased in situations in which the troughs are not readily adaptable to the configuration desired in the facility, or when the troughs, couplings, junctions, downfalls and other equipment do not readily install or easily assemble.\nIn some prior art systems, numerous screws and other attachment hardware are typically utilized to attach the various trough components to one another.\nThere is a need for improved ways to construct and assemble such trough or raceway systems, which is one of the objects of this invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a barrel module of an imaging apparatus, and an imaging apparatus including the barrel module.\n2. Description of the Related Art\nAn apparatus for capturing a digital image stores a still image or a moving image captured of a subject in a digital file. Examples of the apparatus include a digital still camera (DSC), a digital video camera (DVC), and a digital camera installed in a mobile phone.\nGenerally, the apparatus includes a barrel module and a camera body. The barrel module is installed in front of the camera body, and includes an optical lens for inducing and forming the image of the subject, and lens barrels including the optical lens. The lens barrels perform a zoom operation while moving along an optical axis direction. Also, the barrel module includes a focus lens and a focus lens barrel including the focus lens. The focus lens barrel adjusts a focus while moving between an image pickup device and another lens barrel.\nWhen an external force is applied to the barrel module including the plurality of lens barrels, in a direction that the lens barrels are disposed, the lens barrels collide with each other, and thus the barrel module may be damaged."}
{"text": "This invention is directed to swinging bucket rotors and more particularly is directed to a means for reducing the friction between the bucket and the support pins or trunnions when the bucket moves or swings relative to the rotor body.\nPresently it is necessary that the swinging buckets as well as the rotor pins or trunnions that support the buckets be made of strong material to withstand the high pressures and forces that are exerted at the bucket pin interface during the centrifugation. Inherent with these high strength materials is a high friction characteristic.\nDuring the operation of the centrifuge, the swinging bucket moves from a generally vertical orientation to a generally horizontal orientation as the rotor moves from the zero speed to operational speed. Further, when the centrifugation operation is completed and the rotor decelerates, the bucket will move from its horizontal orientation down to its vertical or rest position. When the bucket is moving between its vertical and its horizontal orientations, there is constant relative movement between the bucket and the stationary pins. Because of the high friction characteristic of the materials used, this tends to cause galling at the bucket pin interface. The presence of galling will tend to prevent the bucket from returning to a complete vertical orientation, resulting in possible remixing of the centrifuged sample. Further, this galling can also cause a shortening of the life of the bucket by initiating cracks in the bucket material."}
{"text": "The production of polymers such as polyethylene requires a high purity feedstock of various components, including monomers and co-monomers. In order to offset some of the costs and maximize production, it can be useful to reclaim and/or recycle some feedstock components from an effluent stream resulting from the polymerization reaction. To accomplish this, the reclaimed effluent streams have conventionally either been routed through a purification process or redirected through other redundant processing steps.\nConventional attempts to industrially produce high purity feedstock components has required the operation of numerous distillation columns, compressors (e.g., to achieve the high pressures needed in such conventional processes), refrigeration units (e.g., to achieve cryogenic temperatures) and various other equipment. As such, the equipment and energy costs associated with feedstock purification represent a significant proportion of the total cost for the production of such polymers. Further, the infrastructure required for producing, maintaining, and recycling high purity feedstock represents a significant portion of the associated cost.\nFurther, such conventional attempts to recover feedstock components have not enabled sufficient control parameters to prevent and/or control deleterious plant conditions. The drawbacks of these designs can lead to process delays, increased costs, and/or other inefficiencies. As such, an improved separation system for polymerization reaction effluent streams is needed."}
{"text": "Battery chargers for hand held electronic devices, such as cell phones, portable digital assistants, digital cameras and the like typically include an AC adapter with a wire and wall outlet plug extending from one side of the adapter, and a wire with a coaxial or other connector that can be plugged directly into the device extending from the other side of the adapter. Other chargers may include a charger stand on which the device is placed that establishes an electrical connection with the device, with a wire and wall outlet plug extending from the charger stand. In chargers of the latter type, contacts on the device itself and/or the charger stand, are typically exposed and subject to damage. In addition, for hand-held electronic video game devices, the use of a charger stand typically precludes game operation during the charging process."}
{"text": "This invention relates to a valve block assembly of the kind having multiple inlets and a single outlet. Such assemblies find particular application in protein sequencers.\nAutomated chemistry instrumentation has traditionally used many subsystems to accomplish fluid handling. These subsystems include valving, tubing, tubing connectors, manifolds and fluid reservoirs. Many times, the type of chemistry used in these external sources is sensitive to contamination from internal or external sources so it is critical to design the instrument so that contamination is eliminated or minimized. For example, in the chemistry used in protein sequencers, the so-called Edman degradation, it is extremely important to exclude oxygen from the reaction. The Edman degradation consists of two principal chemical reactions called the coupling reaction and the cleavage reaction. During the coupling reaction, phenylisothiocyanate (PITC) reacts with the protein amino groups in the presence of a base to form a phenylthiocarbamyl (PTC) derivative of the amino-terminal amino acid. The cleavage reaction results in the anilinothiazolinone (ATZ) derivative of the amino terminal amino acid being formed as it is cleaved from the protein chain. The PTC group is extremely sensitive to desulfurization by oxidation so it is imperative that the reaction be performed in the absence of oxygen. If the PTC group is oxidized, the degradation will halt because the ATZ cannot be formed.\nTo overcome the problem of oxygen contamination, it has become common practice to flush the reaction chamber with an inert gas such as nitrogen or argon. While this helps keep the chamber atmosphere oxygen free, it does not prevent oxygen from entering the reaction chamber as a dissolved gas in the reactants themselves. Since the chemicals used in the Edman degradation are also very corrosive, all tubing in the fluid-handling system must be chemically inert. Unfortunately, the tubing, although extremely resistant to attack by the Edman chemicals, is also very porous to the oxygen in the atmosphere and the chemicals become contaminated by the diffusing oxygen.\nPrevious protein sequencers have used tubing to connect remote delivery valves, pressure valves and venting valves to chemical reservoirs with the result that oxygen sensitive chemicals are directly exposed to the diffusing oxygen, not only in the tubing, but also in the reservoirs. There is a direct pathway from the atmosphere through the tubing to the contents of the reservoir from the exposed tubing. To help overcome this problem, it has been a common practice to add reducing agents to all of the chemical reservoirs. While this seems to help significantly, it is not an ideal solution to the problem. The reducing agents contribute spurious components which may interfere with the analysis of the amino acid derivative produced by the sequencer. Also, the antioxidizing effect of the reducing agents is relatively short-term. If the chemicals in the tubing are exposed for a period of days, as in the case of an idle instrument, the chemical must be replaced for the performance of the Edman degradation is greatly compromised.\nAnother source of oxygen contamination is leaking at the junctions between connecting tubing and other components of the fluid-handling system. Previous fluid handling systems have made extensive use of individual components connected through tubing. Every connection represents a potential source of trouble, not only from oxygen contamination but from poor performance due to variable flow rates in a leaking system. Also, since these chemicals may be corrosive or poisonous, a leak represents a hazard to operators and to the instrument itself.\nWittman-Liebold, U.S. Pat. No. 4,008,736 describes a valve arrangement in which a common conduit is formed in the valve block. All delivery valve sites lie on the same surface and are connected by zig-zagging portions of the conduit. While this type of common conduit may be cleaned by flushing, it cannot be claned easily with a wire. It is also difficult to machine.\nGraffunder, U.S. Pat. No. 4,168,724 replaced the slider valves of the '736 patent with diaphragm valves. These valves are closed by fluid pressure in an actuator chamber adjacent to the diaphragm, and opened by evacuating the chamber. The zigzag sections of the common conduit intersect at the surface of the valve block. All delivery valve sites lie on the same surface.\nThis apparatus has several disadvantages. First, it is necessary to provide vacuum and high pressure sources. Second, the diaphragm has a tendency to cold flow into the common conduit at the point of intersection, requiring that it be moved a greater distance in order to open the valve. This in turn places greater demands on the evacuation system, and increases the wear and tear on the diaphragm.\nHunkapiller, U.S. Pat. No. 4,558,845 replaces the zigzag sections of Wittman-Liebold with straight sections which are easier to clean and to machine. However, each valve site is placed on a separate block, and the common conduit is alternately a channel in a block and tubing exposed to the atmosphere.\nHunkapiller also replaced Wittmann's actuation mechanism with a plunger that is spring-biased to a closed position. A solenoid device is used to draw the plunger to an open position. Since Hunkapiller retains the narrow access port of the previously described apparatus, his valve is likewise subject to the problem of membrane \"cold flow\"."}
{"text": "1. Field of the Invention\nThe invention herein relates to tote bags and more particularly to a tote bag that is adapted to secure to the handle of a cart or similar object.\n2. Overview of Prior Art\nA problem that has plagued all of us at one time or another is going to a market and being tasked with the juggling act of hanging onto your wallet, keys and coupons. This is only compounded if you have small children that accompany the shopper. Not only does the task of keeping track of the children increase shopping stress experience but small children commonly sit in the child's seat where many of these items are temporarily stored by the shopper. In an attempt to assist the shopper many shopping bags have been developed, but all fall short of the critical shortfalls presented.\nOne such attempt is disclosed by Posner in U.S. Pat. No. 4,871,100. Here a foldable shopping bag that fits in a shopping cart is shown. Items are placed directly into the bag which is located in the cart. Since the items have to be removed from the cart to be checked and then bagged anyway, the practical function of the invention is limited at best. No disclosure is made to store the shopper's personal items, only the purchased items. To do so is clearly not anticipated and as disclosed does not address the problem as stated.\nA similar theme is disclosed in U.S. Pat. No. 5,531,366 to Strom. Instead of a foldable bag, a collapsible caddie is disclosed that can be folded into a backpack. The function is similar to Posner in that the apparatus fits in the cart and is supported by the upper rim of the basket. Another shortcoming of the disclosed is that the individual compartments minimize the functional capacity of the cart. One of the compartment could be used to store personal effects but the device does not allow for a child to sit in the child's seat while the device is in use. Also the physical size of the carts vary from manufacturer to manufacturer. This device does not appear to anticipate this variance in carts which contributes to the limited function of the device.\nA shopping bag that can be secured to a shopping cart is disclosed by Geeck in U.S. Pat. No. 5,009,516. This is a collapsible bag comprising a large storage bag and a smaller second bag. The combination also has a snap to secure the user's vehicle keys and can be fastened to the inside of a shopping cart. As before, the device is intended to primarily hold purchased items, not personal effects the shopper brings into the store. The device as disclosed, does not present an easily accessible storage for those items to the shopper or a child sitting in the child's seat.\nA shopping bag and tote therefore was disclosed by Caligiuri in U.S. Pat. No. 5,297,872. The disclosure is of a reusable shopping bag, again for purchased items, and a tote bag for storage of the bag when it is not in use. It does not disclose nor anticipate storage of personal items nor any attachment to a shopping cart. In light of this the disclosure does not address the shortcomings as expressed herein.\nIn terms of storage items for the children, Sonders, in U.S. Pat. No. 5,326,300 disclosed a carrying device that is in the form of a child's toy. The toy, such as a teddy bear, includes an internal pocket, suitable for storage of children's items such as a bottle and diapers. The device is clearly not intended to be adapted to a shopping cart nor to the function of serving the shopper and/or a child passenger with storage and transport of personal effects. In that sense it also does not satisfy the need as disclosed.\nAnother child's transport bag was disclosed by Smallwood in U.S. Pat. No. 4,874,340. Here, as previously disclosed, the intended function of the device is to transport items suitable for a child. Toys, books and the like can be stored inside with a toy animal outside to make the device pleasing to young children. The device also includes a pair of straps to enable the device to be used as a backpack. The shortcomings here include that of the device disclosed by Sonders. The device is not intended to be adapted to a shopping cart nor to the function of serving the shopper and/or a child passenger.\nA product from Right Start called the Shopping Cart Friend is also disclosed. This consists of a single flap that attaches to the handle of the cart and is suitable only for that specific child in that toys include rings and such that the child would put into his or her mouth. It would not be desirable to have another child contact this after another just handled it. This necessitates the device be removable and go with the shopper, not allowing a product to stay on the cart permanently. There are also no pockets provided to store items for the child and more importantly the item is intended for use only for the child, there is not another side for use by the shopper.\nA utility bag was disclosed by L,each in U.S. Pat. No. 5,103,514 where the device is comprised of a flexible, rectangular fabric tube with a resilient, flexible material inside. The tube can be folded into a U-shape for transport and carried with the straps which extend outwardly from the sides of the tube. The device is clearly not intended to be associated with a shopping cart nor to the easy, simultaneous accessibility of shopper's personal affects and children's items and as such does not address the disclosed shortcomings in the art.\nOther foldable or collapsible travel pouches are disclosed by Robitille and Franco in U.S. Pat. Nos. 5,244,278 and 4,991,245 respectively. Both disclose a foldable outer case with compartments located therein. Neither shows an anticipated use in conjunction with a shopping cart nor duel simultaneous use for the shopper and a child in the child's seat of the cart.\nA duel cupped coupon holder is disclosed by Anatra in U.S. Pat. No. 5,048,736 with the intended use on a shopping cart handle. The limitations of such a device are far reaching in that the molded one piece container is intended to support only coupons or other light objects. This disclosure does not take into account a child in the seat of the cart. A colorful grouping of individual pieces of paper would greatly attract a small child's attention. The disclosure does not provide for security for valuables such as money or other personal items from a child in the seat nor from a thief posing as a shopper.\nA bicycle brief case is disclosed by Tiffany in U.S. Pat. No. 5,031,807 in which a saddle bag construction of a briefcase is fastened to the frame of a bicycle. This would be less than adequate for use on a shopping cart because of the intended function of the fasteners being at the bottom of the flaps and as before the pockets open at the top of the device. This allows access to anyone who can get their hand into the opening of the bag, even a child in the seat of the cart."}
{"text": "A significant number of adults have had an episode of back pain or suffer chronic back pain emanating from a region of the spinal column. A number of spinal disorders are caused by traumatic spinal injuries, disease processes, aging processes, and congenital abnormalities that cause pain, reduce the flexibility of the spine, decrease the load bearing capability of the spine, shorten the length of the spine, and/or distort the normal curvature of the spine. Many people suffering from back pain resort to surgical intervention to alleviate their pain.\nDisc degeneration can contribute to back pain. With age, the nucleus pulposus of the intervertebral discs tends to become less fluid and more viscous. Dehydration of the intervertebral disc and other degenerative effects can cause severe pain. Annular fissures also may be associated with a herniation or rupture of the annulus causing the nucleus to bulge outward or extrude out through the fissure and impinge upon the spinal column or nerves (i.e. a “ruptured” or “slipped” disc).\nIn addition to spinal deformities that can occur over several motion segments, spondylolisthesis (i.e. forward displacement of one vertebra over another, usually in the lumbar or cervical spine) is associated with significant axial and/or radicular pain. Patients who suffer from such conditions can experience diminished ability to bear loads, loss of mobility, extreme and debilitating pain, and oftentimes suffer neurological deficit in nerve function.\nFailure of conservative therapies to treat spinal pain such as for example bed rest, pain and muscle relaxant medication, physical therapy or steroid injection often urges patients to seek spinal surgical intervention. Many surgical techniques, instruments and spinal disc implants have been described that are intended to provide less invasive, percutaneous, or minimally-invasive access to a degenerated intervertebral spinal disc. Instruments are introduced through the annulus for performing a discectomy and implanting bone growth materials or biomaterials or spinal disc implants within the annulus. One or more annular incisions are made into the disc to receive spinal disc implants or bone growth material to promote fusion, or to receive a pre-formed, artificial, functional disc replacement implant.\nExtensive perineural dissection and bone preparation can be necessary for some of these techniques. In addition, the disruption of annular or periannular structures can result in loss of stability or nerve injury. As a result, the spinal column can be further weakened and/or result in surgery-induced pain syndromes.\nVarious posterior interbody approaches such as posterior lumbar interbody fusion (PLIF) and transforaminal lumbar interbody fusion (TLIF) necessitate placing an interbody device through a relatively small annular aperture from a posterior access dissection path. The desire to restore disc space height and create lordosis through anterior distraction is in conflict with the application of a fixed height interbody spacer through a limited posterior approach. In addition, the larger the “foot-print” of the interbody spacer, the less likely it will subside. This is not only due to reduced endplate pressure but also better load bearing bone near the perimeter of the endplate (versus the central region)."}
{"text": "Conventionally, in application software running on a computer, there are cases in which input of parameters for command instructions for processing, input or changes of properties, or the like are accepted. As a result, pop-up menus are used in many cases as one type of user interface.\nFor pop-up menus, a screen is displayed in response to, for example, a right-click of a mouse operated by a user, and by selection of a processing item from a displayed list, a corresponding process is executed. At that time, in a case where parameter input is required, a dialog box provided for this purpose is displayed, and input is accepted while execution of processing is triggered by pressing of the “run” button or the like provided in the dialog box.\nHere, in application software, the more functions provided, the more processing items are displayed in pop-up menus, so it can also be considered that, for a user, an amount of time to search for a desired processing item increases.\nTherefore, a method is proposed for improving a user interface with this kind of possibility of reducing convenience. For example, in Patent Document 1, a method is disclosed in which, by enabling effective display of a shortcut menu (pop-up menu) composed of a plurality of (many) levels, operability of the shortcut menu is increased, and it becomes possible for a user to carry out the user's desired processes more quickly.\n[Patent Document 1] Japanese Unexamined Patent Application Publication No. 2006-59155"}
{"text": "The process known as electrophoresis involves the migration of charged molecules through a suitable retarding medium under the influence of an electric field. Generally, the compounds of higher molecular weight migrate at a slower rate through the medium than do the compounds of lower molecular weight. Devices have been provided previously for carrying out electrophoresis. An example of such a device is U.S. Pat. No. 4,415,418 in which a tray is provided with a raised platform at the center. Removable partitions are placed in the tray at opposite ends of the platform, and a conventional electrophoresis gel is poured over the platform to form a thin layer. When the gel has cooled, the partitions are removed. A comb is provided to form wells across the surface of the gel. Substances that are to be subjected to electrophoresis are delivered into each of the wells, and the tray is at least partially filled with an electrolyte buffer. Electrodes are positioned at each end of the tray and a sufficient voltage difference is applied to the electrodes to cause migration of the molecules of the substance in the wells across the length of the gel, separated according to their molecular weight. After electrophoresis, the gel is removed from the original casting tray, and placed in a dish containing depurination solution. Approximately thirty minutes later this solution is poured out by tipping the dish toward one edge while the gel is held with the fingers. It is important to use great care during this procedure to prevent the gel from breaking because there is no gel support structure and subsequent processing is possible only with an integral gel. A denaturation solution is then added to the dish and incubation is continued for approximately thirty minutes. Again, the solution is carefully poured off. Then neutralization buffer is added and incubation is continued for thirty additional minutes.\nIn accordance with conventional techniques, transfer of the nucleic acids is accomplished by placing a piece of filter paper, which is as wide as and longer than the gel, on a platform which is suspended above a solution of 10.times. saturated saline citrate buffer (SSC). The ends of the filter paper are long enough to hang off the ends of the platform and dip into the 10.times. SSC. Thus, the filter paper acts as a wick to absorb the SSC solution. The gel is removed from the dish and placed on top of the filter paper saturated with 10.times. SSC. Next, a piece of membrane filter paper which is the same size as the gel is saturated with 10.times. SSC and placed on top of the gel. The nucleic acids are eventually bound to the membrane filter paper. Another piece of saturated filter paper, the same size as the gel, is placed on top of the membrane. The entire layered unit is then smoothed to remove any air bubbles that may exist between the gel and the filter paper. Finally, a stack of paper towels, the same size as the gel, is positioned on top of the layered unit.\nOver a period of about 12 to 16 hours, the 10.times. SSC solution is drawn up through the gel by capillary action and the nucleic acids are transferred out of the gel into the membrane above. The paper towels absorb the excess buffer and provide the force for capillary action. At the end of the transfer period, the entire layered unit is disassembled and the membrane is removed for hybridization. This technique is described in an article by E. Southern, \"Detection of Specific Sequences Among DNA Fragments Separated by Gel Electrophoresis,\" J. Mol. Biol., 98:503 (1975).\nAlthough the trays such as the one described in U.S. Pat. No. 4,415,418 are convenient for carrying out electrophoresis, they are not suitable for situations where a large number of samples must be tested in a relatively short period of time.\nTherefore, the prior art uses a tedious multi-step, multi-apparatus process for preparing nucleic acid fragments for subsequent hybridization. Four steps are generally undertaken to achieve preparation of the sample for hybridization. Electrophoresis was previously described. Depurination removes purine bases from nucleic acids. Denaturation involves separating the strands of nucleic acids and breaks down the depurinated nucleic acids into suitable size to allow eventual transfer of the fragments out of the gel. Transfer involves allowing the fragments to go out of the gel onto the porous membrane."}
{"text": "Applications that are distributed via an application marketplace are often installed on multiple devices by a user. A developer of an application may wish to synchronize state information across the multiple installations across various devices of a particular application. Accordingly, developers can employ one or more software libraries for the purpose of generating state information of an application as the application executes, which can be synchronized across the various installations of the application. The state information generated by various devices associated with the user may create synchronization conflicts."}
{"text": "Multi-lumen catheters are used for the purpose of creating two or more separate fluid pathways, such as in hemodialysis applications. A primary goal of hemodialysis access is to provide a reliable and effective means of dialysis, which means that a sufficient volume of blood over a period of time must be removed from and returned to the patient. Because the contaminated and cleansed blood must be kept separate for an effective dialysis procedure, a dual lumen catheter is generally used. These dual lumen catheters are usually configured so that there is a shorter lumen that aspirates blood from a blood vessel of a patient to a dialysis machine where it is processed for the removal of toxins, and a longer lumen that infuses the purified blood to the patient. The shorter lumen utilized for aspiration is generally referred to as the “arterial lumen,” while the longer lumen utilized for infusion is generally referred to as the “venous lumen.” The reason for the different lengths is to minimize co-mingling of aspirated and infused blood.\nThe primary problems occurring in dual lumen dialysis catheters are blood clotting (thrombosis) and fibrin (the protein formed during normal blood clotting that is the essence of the clot) sheath formation. Thrombus and fibrin sheath formation can occlude distal tips of the dialysis catheter lumens, resulting in loss of catheter function when such an occlusion prevents blood flow. This typically occurs initially in the arterial lumen used for aspiration of blood from a patient. A secondary problem relates to the arterial lumen “sucking” against the vessel wall, in which the arterial lumen openings become fully occluded by the patient's vasculature.\nTo specifically address these problems, a new type of dialysis access catheter has been designed that utilizes independent “free floating” distal tip sections that separate at a distal end of the catheter to theoretically reduce the likelihood of potential occlusion and “sucking” during dialysis treatment. U.S. Pat. No. 6,001,079 to Pourchez and U.S. Pat. Nos. 5,947,953 and 6,190,349 to Ash et al., all incorporated by reference herein, are directed to said new type of catheter, hereinafter referred to as a “split-tip catheter.”\nCatheters used for hemodialysis are generally inserted according to standard technique. This includes identifying the location of the central veins in which the catheter is to be inserted, utilizing the Seldinger technique to cannulate the central vein with a guidewire, and passing various instruments over the guidewire until the catheter is in place. In the case of a small diameter catheter made of relatively stiff material (i.e., a percutaneous catheter for acute access), the catheter itself can be passed over the guidewire and into the central vein. In the case of a larger diameter catheter, the passage to the central vein must be widened through the use of dilator(s), after which the catheter can be passed over the guidewire. Another possibility available to the physician is to use an introducer sheath, which is passed over the guidewire (after dilation or simultaneously therewith) forming a path for the catheter to access the central vein.\nDue to the configuration of the split-tip catheter, standard insertion techniques must be modified. With respect to non-sheath insertions, the separate distal tip sections must each be passed over the guidewire to ensure their arrival at the targeted central vein without attendant problems. Suggested means for accomplishing this (in a dual lumen catheter) is to pass the back end of the guide wire through the distal end hole of one tip section and out a side hole thereof into the distal end hole of the other tip section (Patel et al., J. Vasc. Interv. Radiol. (2001), vol. 12, pages 376-378). However, such a method can become problematic depending on the placement of the side holes and, in fact, may be impossible in the event that no side holes are provided. With respect to sheathed insertions, the separate distal tip sections may tend to flare outward and catch or “snag” on the inner wall of the delivery sheath thereby preventing smooth delivery of the catheter to the central vein.\nThere is a need for improvements to split-tip catheters and new configurations thereof, which enhance the design and further alleviate concerns with regard to typical problems encountered in hemodialysis catheters. There is also a need for new configurations and new methods for insertion of split-tip catheters to overcome problems associated therewith"}
{"text": "This invention relates to an improved load cell of the type used to measure tension or compression loads and more particularly relates to a washer type load cell adapted to be used with gap-spanning sensor elements.\nLoad cells have many applications and are used in the prior art in electronic weighing systems for monitoring force and pressure and so on. As such, these cells employ a suitable sensor element such as piezoelectric, reluctance, capacitive or piezoresistive element which is positioned in a housing to monitor the force or pressure exerted on the housing when the load cell is positioned in a suitable environment.\nThe most popular types of load cells employ a strain gage element which may be an unbonded or bonded wire gage or a foil or a bonded semiconductor element. These gages provide higher sensitivity while further eliminating mechanical linkages and moving parts. They also enable one to achieve good temperature specifications and serve to reproduce load deflection in the elastic member or housing with good linearity.\nA major concern in the design of the load cell resides in the housing arrangement. In this manner, the housing serves to couple or transmit the force or pressure to be monitored by the cell to the sensor associated with the cell. The housing would therefore be capable of transmitting such forces to the sensor in a linear manner without providing undue distortions and improper bending movements when the monitored load is not properly aligned with the axis of the load cell. Some load cells are required to have high overall stiffness. Thus, the load cell housing should provide a high overall stiffness while still allowing sufficient strain to be monitored by the associated sensor element.\nIt is therefore an object of the present invention to provide a load cell transducer capable of accurately monitoring an external force employing a housing which is simple and relatively inexpensive to fabricate; and which provides a relatively stiff structure, but also allows for adequate strain levels to be transmitted to the sensor element or elements."}
{"text": "A NAND type flash memory with a three-dimensional structure manufactured using a BiCS (Bit Cost Scalable Flash memory: BiCS) manufacturing technique is referred to in the art as BiCS memory."}
{"text": "1. Field of the Invention\nThe present invention relates data communications and, more particularly, to data communications using multicarrier modulation.\n2. Description of the Related Art\nBi-directional digital data transmission systems are presently being developed for high-speed data communication. One standard for high-speed data communications over twisted-pair phone lines that has developed is known as Asymmetric Digital Subscriber Lines (ADSL). Another standard for high-speed data communications over twisted-pair phone lines that is presently proposed is known as Very High Speed Digital Subscriber Lines (VDSL).\nThe Alliance For Telecommunications Information Solutions (ATIS), which is a group accredited by the ANSI (American National Standard Institute) Standard Group, has finalized a discrete multi tone based approach for the transmission of digital data over ADSL. The standard is intended primarily for transmitting video data and fast Internet access over ordinary telephone lines, although it may be used in a variety of other applications as well. The North American Standard is referred to as the ANSI T1.413 ADSL Standard (hereinafter ADSL standard). Transmission rates under the ADSL standard are intended to facilitate the transmission of information at rates of up to 8 million bits per second (Mbits/s) over twisted-pair phone lines. The standardized system defines the use of a discrete multi tone (DMT) system that uses 256 xe2x80x9ctonesxe2x80x9d or xe2x80x9csubchannelsxe2x80x9d that are each 4.3125 kHz wide in the forward (downstream) direction. In the context of a phone system, the downstream direction is defined as transmissions from the central office (typically owned by the telephone company) to a remote location that may be an end-user (i.e., a residence or business user). In other systems, the number of tones used may be widely varied. However when modulation is performed efficiently using an inverse fast Fourier transform (IFFT), typical values for the number of available sub-channels (tones) are integer powers of two, as for example, 128, 256, 512, 1024 or 2048 sub-channels.\nThe ADSL standard also defines the use of a reverse signal at a data rate in the range of 16 to 800 Kbit/s. The reverse signal corresponds to transmission in an upstream direction, as for example, from the remote location to the central office. Thus, the term ADSL comes from the fact that the data transmission rate is substantially higher in the downstream direction than in the upstream direction. This is particularly useful in systems that are intended to transmit video programming or video conferencing information to a remote location over telephone lines.\nBecause both downstream and upstream signals travel on the same pair of wires (that is, they are duplexed) they must be separated from each other in some way. The method of duplexing used in the ADSL standard is Frequency Division Duplexing (FDD) or echo canceling. In frequency division duplexed systems, the upstream and downstream signals occupy different frequency bands and are separated at the transmitters and receivers by filters. In echo cancel systems, the upstream and downstream signals occupy the same frequency bands and are separated by signal processing.\nANSI is producing another standard for subscriber line based transmission system, which is referred to as the VDSL standard. The VDSL standard is intended to facilitate transmission rates of at least about 6 Mbit/s and up to about 52 Mbit/s or greater in the downstream direction. To achieve these rates, the transmission distance over twisted-pair phone lines must generally be shorter than the lengths permitted using ADSL. Simultaneously, the Digital, Audio and Video Council (DAVIC) is working on a similar system, which is referred to as Fiber To The Curb (FTTC). The transmission medium from the xe2x80x9ccurbxe2x80x9d to the customer is standard unshielded twisted-pair (UTP) telephone lines.\nA number of modulation schemes have been proposed for use in the VDSL and FTTC standards (hereinafter VDSL/FTTC). For example, some of the possible VDSL/FTTC modulation schemes include multi-carrier transmission schemes such as Discrete Multi-Tone modulation (DMT) or Discrete Wavelet Multi-Tone modulation (DWMT), as well as single carrier transmission schemes such as Quadrature Amplitude Modulation (QAM), Carrierless Amplitude and Phase modulation (CAP), Quadrature Phase Shift Keying (QPSK), or vestigial sideband modulation.\nMost of the proposed VDSL/FTTC modulation schemes utilize frequency division duplexing of the upstream and downstream signals. One particular proposed VDSL/FTTC modulation scheme uses periodic synchronized upstream and downstream communication periods that do not overlap with one another. That is, the upstream and downstream communication periods for all of the wires that share a binder are synchronized. When the synchronized time division duplexed approach is used with DMT it is referred to as synchronized DMT (SDMT). With this arrangement, all the very high speed transmissions within the same binder are synchronized and time division duplexed such that downstream communications are not transmitted at times that overlap with the transmission of upstream communications. This is also referred to as a (i.e. xe2x80x9cping pongxe2x80x9d) based data transmission scheme. Quiet periods, during which no data is transmitted in either direction, separate the upstream and downstream communication periods.\nA common feature of the above-mentioned transmission systems is that twisted-pair phone lines are used as at least a part of the transmission medium that connects a central office (e.g., telephone company) to users (e.g., residence or business). It is difficult to avoid twisted-pair wiring from all parts of the interconnecting transmission medium. Even though fiber optics may be available from a central office to the curb near a user\"\"s residence, twisted-pair phone lines are used to bring in the signals from the curb into the user\"\"s home or business.\nThe twisted-pair phone lines are grouped in a binder. While the twisted-pair phone lines are within the binder, the binder provides reasonably good protection against external electromagnetic interference. However, within the binder, the twisted-pair phone lines induce electromagnetic interference on each other. This type of electromagnetic interference is generally known as crosstalk interference which includes near-end crosstalk (NEXT) interference and far-end crosstalk (FAR) interference. As the frequency of transmission increases, the crosstalk interference becomes substantial. As a result, the data signals being transmitted over the twisted-pair phone lines at high speeds can be significantly degraded by the crosstalk interference caused by other twisted-pair phone lines in the binder. As the speed of the data transmission increases, the problem worsens.\nMulticarrier modulation has been receiving a large amount of attention due to the high data transmission rates it offers. FIG. 1A is a block diagram of a conventional transmitter 100 for a multicarrier modulation system. The transmitter 100 receives data signals to be transmitted at a buffer 102. The data signals are then supplied from the buffer 102 to a forward error correction (FEC) unit 104. The FEC unit 104 compensates for errors that are due to crosstalk noise, impulse noise, channel distortion, etc. The signals output by the FEC unit 104 are supplied to a data symbol encoder 106. The data symbol encoder 106 operates to encode the signals for a plurality of frequency tones associated with the multicarrier modulation. In assigning the data, or bits of the data, to each of the frequency tones, the data symbol encoder 106 utilizes data stored in a transmit bit allocation table 108 and a transmit energy allocation table 110. The transmit bit allocation table 108 includes an integer value for each of the carriers (frequency tones) of the multicarrier modulation. The integer value indicates the number of bits that are to be allocated to the particular frequency tone. The value stored in the transmit energy allocation table 110 is used to effectively provide fractional number of bits of resolution via different allocation of energy levels to the frequency tones of the multicarrier modulation. In any case, after the data symbol encoder 106 has encoded the data onto each of the frequency tones, an Inverse Fast Fourier Transform (IFFT) unit 112 modulates the frequency domain data supplied by the data symbol encoder 106 and produces time domain signals to be transmitted. The time domain signals are then supplied to a digital-to-analog converter (DAC) 114 where the digital signals are converted to analog. Thereafter, the digital signals are transmitted over a channel to one or more remote receivers.\nFIG. 1B is a block diagram of a remote receiver 150 for a conventional multicarrier modulation system. The remote receiver 150 receives analog signals that have been transmitted over a channel by a transmitter. The received analog signals are supplied to an analog-to-digital converter (ADC) 152. The ADC 152 converts the received analog signals to digital signals. The digital signals are then supplied to a Fast Fourier Transform (FFT) unit 154 that demodulates the digital signals while converting the digital signals from a time domain to a frequency domain. The demodulated digital signals are then supplied to a frequency domain equalizer (FEQ) unit 156. The FEQ unit 156 performs an equalization on the digital signals so the attenuation and phase are equalized over the various frequency tones. Then, a data symbol decoder 158 receives the equalized digital signals. The data symbol decoder 158 operates to decode the equalized digital signals to recover the data, or bits of data, transmitted on each of the carriers (frequency tones). In decoding the equalized digital signals, the data symbol decoder 158 needs access to the bit allocation information and the energy allocation information that were used to transmit the data. Hence, the data symbol decoder 158 is coupled to a received bit allocation table 162 and a received energy allocation table 160 which respectively store the bit allocation information and the energy allocation information that were used to transmit the data. The data obtained from each of the frequency tones is then forwarded to the forward error correction (FEC) unit 164. The FEC unit 164 performs error correction of the data to produce corrected data. The corrected data is then stored in a buffer 166. Thereafter, the data may be retrieved from the buffer 166 and further processed by the receiver 150. Alternatively, the received energy allocation table 160 could be supplied to and utilized by the FEQ unit 166.\nOne problem with the conventional design of transmitters and receivers of multicarrier modulation systems such as illustrated in FIGS. 1A and 1B is that only a single bit allocation is provided for transmission or reception of data symbols. In particular, the transmitter 108 has a single set of bit allocation information stored in the transmit bit allocation table 108 and the receiver 200 has a corresponding single set of bit allocation information stored in the receive bit allocation table 212. Although the bit allocation table is changeable, the processing time to update or change bit allocations is relatively slow and typically requires some sort of training process. With only a single bit allocation available to the multicarrier modulation system, the multicarrier modulation system is unable to rapidly alter its bit allocations for symbols being transmitted and received. In other words, during transmission or reception of data, the bit allocations are fixed, and thus, all symbols being transmitted and received must use the same bit allocations.\nThus, there is a need for improved transmitters and receivers of multicarrier modulation systems that are able to support multiple bit allocations so that multicarrier modulation systems are able to rapidly alter their bit allocations.\nBroadly speaking, the invention is a method and apparatus for supporting multiple bit allocations in a multicarrier modulation system so that all symbols being transmitted or received can make use of different bit allocations. By supporting the multiple bit allocations, the multicarrier modulation system is able to support bit allocation on a superframe basis. The invention also pertains to selection and alignment of superframe formats to improve system performance. The invention is suitable for use in data transmission systems in which transmissions use a frame structure. The invention is also well suited for data transmission systems involving different transmission schemes where multiple bit allocations are helpful to reduce crosstalk interference.\nThe invention can be implemented in numerous ways, including as an apparatus, system, method, or computer readable media. Several embodiments of the invention are discussed below.\nAs a transmitter for a data transmission system using multicarrier modulation, one embodiment of the invention includes: a superframe bit allocation table, a data symbol encoder, a multicarrier modulation unit, and a digital-to-analog converter. The superframe bit allocation table stores superframe bit allocation information including separate bit allocation information for a plurality of frames of a superframe. The data symbol encoder receives digital data to be transmitted and encodes bits associated with the digital data to frequency tones of a frame based on the superframe bit allocation information associated with the frame stored in said superframe bit allocation table. The multicarrier modulation unit modulates the encoded bits on the frequency tones of a frame to produce modulated signals. The digital-to-analog converter converts the modulated signals to analog signals.\nAs an apparatus for recovering data transmitted by a transmitter, an embodiment of the invention includes: an analog-to-digital converter, a demodulator, a superframe bit allocation table, and a data symbol decoder. The analog-to-digital converter receives transmitted analog signals and produces digital signals therefrom, the transmitted analog signals being time domain signals representing data transmitted. The demodulator receives the digital signals and demodulates the digital signals to produce digital frequency domain data. The superframe bit allocation table stores superframe bit allocation information including separate bit allocation information for a plurality of frames of a superframe. The data symbol decoder operates to decode bits associated with the digital frequency domain data from frequency tones of a frame based on the superframe bit allocation information associated with the frame stored in said superframe bit allocation table.\nAs a method for allocating bits to symbols of a superframe for transmission of data in a data transmission system using multicarrier modulation, an embodiment of the invention includes the operations of: receiving a service request for data transmission; determining a number of bits required to support the service request; obtaining performance indicia for a plurality of the symbols in a superframe; and allocating the determined number of bits to a plurality of symbols in the superframe based on the performance indicia.\nAs a method for determining an alignment for a superframe used to transmit data in a data transmission system using multicarrier modulation, an embodiment of the invention includes the operations of: (a) receiving a service request for data transmission; (b) selecting a superframe format based on the service request; (c) selecting a proposed alignment of the selected superframe format; (d) allocating bits to frequency tones of the selected superframe format; (e) determining a performance measure for the selected superframe format with the allocation of bits; (f) repeating operations (c)-(e) for at least one other proposed alignment; (g) choosing the one of the proposed alignments for the superframe format in accordance with the determined performance measures.\nAs a method for allocating bits to symbols of a superframe for transmission of data in a data transmission system using multicarrier modulation, an embodiment of the invention includes the operations of: (a) receiving a service request for data transmission; (b) selecting a superframe format based on the service request; (c) determining an alignment of the selected superframe format; (d) allocating bits to frequency tones of the selected superframe format having the alignment; (e) determining a performance measure for the selected superframe format with the allocation of bits; (f) repeating operations (b)-(e) for at least one other superframe format; (g) choosing the superframe format in accordance with the determined performance measures."}
{"text": "1. Field of the Invention\nThe present invention relates to a casing structure for electronic equipment that has the effect of shielding electromagnetic waves. More particularly, the present invention relates to a casing structure for electronic equipment that has the effect of shielding electromagnetic waves by means of electrical grounding between a case member and an electronic equipment support member that serves to support electronic boards and the like when the electronic equipment support member is mounted in the case member\n2. Background Information\nDue to improvements in electronics technology, the use of electronic equipment, consumer electronics, communications equipment, and the like are becoming increasingly common in our daily lives. However, these devices continually generate electromagnetic waves when electrical power is supplied to them. In particular, the increase in the speed of the process and the colorization in recent image forming devices (e.g., digital copying machines, printers, and the like) is one cause of many of the electromagnetic waves generated by various devices arranged inside an image forming device. The electromagnetic waves generated by these devices are capable of having an adverse impact on other circuits in the devices and on other external electronic equipment, and thus there is a need to prevent the exterior leakage of electromagnetic waves to the greatest degree possible.\nControl at a high clock frequency is needed in order to increase the speed of the process in an image forming device. When this occurs, strong electromagnetic waves may be generated from the control board or the like that generates the high clock frequency. For example, with a digital copying machine, multiple overlap of signal forms that are synchronized with the clock frequency will occur due to the relationship between high speed original document scanning and image formation when these are simultaneously performed, and thus will cause even stronger electromagnetic radiation. In addition, with recent color copying machines and printers, electromagnetic waves in an even higher frequency band may be generated at a high intensity due to the overlap of the clock frequency of a charge coupled device control board that is color separated into at least four colors, the increase in speed, and the improvement in image quality (high density pixels). Because of these issues, countermeasures against the strong electromagnetic waves that leak from image forming devices are in urgent demand. For example, electromagnetic wave shielding for the leakage of electromagnetic waves in the 30 to 230 MHz range is now being studied by groups such as the Voluntary Control Council of Japan For Interference By Information Technology Equipment (VCCI). However, the countermeasures for the electromagnetic wave shielding must respond to recent demands for lower cost equipment, and thus it is desirable that the electromagnetic wave shielding maintain superior electromagnetic wave shielding effects while performing the shielding with the most inexpensive means possible.\nVarious types of electromagnetic shielding for electronic equipment have been tried in the past. One example of these is a conductive flexible member that is attached to an image forming device as a grounding member (see, for example, Japanese unexamined patent application publication H5-307282). With Japanese patent application publication H5-307282, one end of the conductive flexible member is attached to the main portion of the electronic equipment with a screw, and when the electronic equipment is grounded, the plate spring effect of the flexible member will be used to press the other end of the flexible member to be in contact with another member.\nIn another example, an elastic conductive member is used as an electromagnetic shielding material, and attached to the contact portions between the support members and frame of electronic equipment to prevent the leakage of electromagnetic waves (see, for example, Japanese unexamined patent application publication H6-291485). With Japanese unexamined patent application publication H6-291485, a conductive sponge is adhered to the contact portions between the support members and case member with a double-sided adhesive tape, the portions of the support members and the case member that are out of range of the adhesive area of the double-sided adhesive tape are brought into contact with each other, and the electrical contact of both is thereby achieved.\nIn another example, an image reading device in which the end portions of a board are bent, and contact pressure is produced on the decorative case of a device during board mounting in order to perform grounding (see, for example, Japanese unexamined patent application publication H6-105080).\nIn addition, other examples include an image forming device in which 2 units are brought into contact via a conductive elastic member (see, for example, Japanese unexamined patent application publication H8-190237), and a structure in which a shielding cover covers the upper portion of an electromagnetic noise generation source (see, for example, Japanese unexamined patent application publication 2002-237915).\nHowever, with the grounding configuration of the image forming device that is disclosed in Japanese unexamined patent application publication H5-307282, it will be necessary to newly provide a grounding member for electromagnetic shielding in the image forming device, and it will also be necessary to fix the grounding member in the image forming device with a large number of screws or the like. Because of this, the number of parts and steps needed during the assembly of the grounding structure of the image forming device disclosed in this reference will increase, the assembly efficiency of the image forming device will worsen, and manufacturing costs will increase.\nWith the electromagnetic wave prevention device disclosed in Japanese unexamined patent application publication H6-291485, double-sided adhesive tape is employed to fix the conductive sponges, but even with this construction, a new and separate member for electromagnetic wave shielding is fixed with a fixing means (double-sided tape), and thus the previously noted problems of worsened assembly efficiency due to the increase in the number of parts and assembly steps, and the increase in manufacturing costs, will not be solved.\nWith the image reading device disclosed in Japanese unexamined patent application publication H6-105080, if there is an increase in emission waves from the board, the equipment inside the device will malfunction due to this increase, the ends of the board will become antennas, and electromagnetic waves will be easily emitted to the exterior of the device. Furthermore, with Japanese unexamined patent application publication H6-105080, although electrical connectivity is achieved by bending the ends of the board, the opposite sides of the bend portions are fixed with screws, and thus this construction cannot contribute to the elimination of parts and a simplified assembly.\nFurthermore, even with the devices of Japanese unexamined patent application publication H8-190237 and Japanese unexamined patent application publication 2002-237915, there will be problems with the ease of assembly, the cost aspect, and the workability thereof, and thus there is a great deal of room for improvement.\nThus, an object of the present invention is to take these problems into consideration, and provide a casing structure for electronic equipment in which the assembly of the electronic equipment is simplified while providing both a grounding effect and an electromagnetic wave shielding effect, and in which the low cost manufacture of the electronic equipment can be achieved."}
{"text": "The present invention relates to a proteome analysis method and a proteome analysis system, which identifies proteins by using a mass spectrometer, in particular, which has the high accuracy and the high throughput, by using predicted results of bond cleavage parts by a molecular simulation method such as a molecular orbital (MO) method."}
{"text": "Silicon carbide (SiC), also known as carborundum, is a compound of silicon and carbon with chemical formula SiC. SiC may be produced in many forms, including but not limited to bulk ceramic materials, grains and filaments. Bulk or granular SiC has numerous uses including use as abrasive or cutting tools, a structural ceramic, in electronic circuit components and heating elements. SiC can also be produced in a whisker form. As is customary in the SiC production industry a “whisker” is defined as a SiC particle having a high aspect ratio of length to diameter. SiC whiskers of various sizes and manufactured by various techniques are often used to reinforce or toughen other materials.\nThe earliest SiC whiskers produced were fabricated from rice hulks (husks) and intense pressure in a SiO2 atmosphere. Later SiC whiskers were produced using petrochemical processes, and most recently production has relied on direct fluorination at high temperatures. A promising approach uses a SiO2—CH4—Na3AlF6 process, for example. Some attention has recently been focused on formation of molybdenum silicide-SiC powders produced in a similar manner.\nThus, known processes use either: (i) very high temperature for reacting constituents into the desired SiC, (ii) expensive and difficult to handle petrochemical derivatives (such as benzoxazine, toluene) or (iii) expensive precursors (e.g., B4C, “boron carbide”). For example, typical known processes occur at greater than 1,000° C., up to 2,200° C. which becomes problematic because higher temperatures demand different furnaces/processes with more critical control-related problems, for example, feed speed, time controlling in heat zone, gas discharge, heat discharge, and power quality monitoring.\nThe embodiments disclosed herein are directed toward overcoming one or more of the problems discussed above."}
{"text": "The present invention relates to a 3-D data storing method, and more particularly to a 3-D data storing method capable of accessing the data of eight cubical vertices belonging to a cube during a memory access cycle.\nThe term 3-D data refers to data groups which are arranged in a mesh structure. This arrangement is chiefly used for color correction in a color television. For interpolating a color value of a point around a known point in the color space expressed by the three primary colors, i.e. red, green and blue, the color values of the known point and eight cubical vertices surrounding the known point should be known.\nAccording to a conventional method, the 3-D data is sequentially stored in a single memory, thus requiring eight memory accessing operations for reading the color values of the eight cubical vertices. That is, to acquire the data of the cubical vertices of a cube, eight memory accessing operations are required. Such a conventional method limits the operating speed of a system for processing 3-D data, thus making real-time processing difficult."}
{"text": "1. Field of the Invention\nThe present invention relates to peripheral device identification and interrupt management. More specifically, the present invention pertains to a method and system that enables device identification to be performed without a rigid time limit within which an initial interrupt from the device needs to be serviced.\n2. Related Art\nComputer systems and other electronic devices have become integral tools that are used extensively to perform a wide variety of useful operations in modern society. Applications of computer systems can be found in virtually all fields and disciplines, including but not limited to business, industry, scientific research, education and entertainment. For instance, computer systems are used to analyze financial data, to control industrial machinery, to model chemical molecules, to deliver classroom presentations and to generate special effects for movies. Moreover, computer systems along with other electronic devices are finding new applications as a result of advances in hardware technology and rapid development in software technology. The growing affordability of computer systems and electronic devices together with the abundance of useful new applications have fueled strong demand for such systems and devices.\nIncluded within this broad category of computers and electronic devices is the personal digital assistant (commonly referred to as a PDA). Specifically, as the components required to build a computer system have been greatly reduced in size, new categories of computer systems have emerged. One of these new categories of computer systems is the PDA. A PDA is a portable computer system which is small enough to be held conveniently and comfortably in the hand of its user. In particular, a popular size for the typical PDA approximates the size of a palm.\nThe PDA is usually a battery-powered device that is typically used as an electronic organizer having the capability to store and display a wide range of information which can include numerous addresses and telephone numbers of business and personal acquaintances, financial information, daily appointments, along with various other personal information. As such, the PDA is able to consolidate a wide variety of information and make the information easily accessible by its user. Therefore, PDAs are very useful and have gained wide popularity.\nSince PDAs are such small devices, full-sized keyboards are generally not efficient input devices. For instance, PDAs using keyboards usually have keyboard devices that are so small that a user typically cannot touch-type on them. Many PDAs thus employ a stylus and a digitizer pad as an input system. The stylus and digitizer pad work well for PDAs because the arrangement allows a user to hold the PDA in one hand while writing with the stylus onto the digitizer pad with the other hand. A small on-screen keyboard image can also be used for data entry. By eliminating the need for a keyboard, PDAs are very portable and can be carried along by their users wherever they go, even when on extended travel. On the other hand, due to their small size, PDAs usually have a modest set of built-in functions and it is often beneficial to couple peripheral devices therewith to expand functionality. The greater the variety of peripheral devices a PDA is capable of operating with, the broader the functionality it can achieve.\nMoreover, the management of communications port contention and power consumption affect the ease of use and the battery life of PDAs and other hand-held computers. A typical PDA includes limited communications port resources. For example, the Palm V personal organizer, manufactured by Palm Computing, Inc. of Santa Clara, Calif., has a single communications port available to interface with peripherals. When the communications port is exclusively occupied by a peripheral, the communications port is unavailable for other uses. Thus, a PDA with an active keyboard occupying the communications port, for instance, may be unable to perform an infrared (IR) synchronization process with another device.\nAnother ease of use issue is the degree of user intervention required to manage a communications port and identify peripheral devices coupled thereto. Requiring a user to manually open the communications port and identify a peripheral device, such as through the graphical user interface (GUI) or the buttons of the PDA or hand-held computer, is generally less preferable to automatically opening the communications port and identifying the peripheral device when the PDA receives information from a peripheral device upon coupling. Further, requiring a user to explicitly close the communications port is also generally less desirable than automatically closing the communications port after the peripheral device and the PDA no longer exchange data. In one existing solution, the PDA will automatically close a communications channel, which was opened to work with a keyboard peripheral, when the PDA is prompted with a signal initiating a synchronization process, such as when the PDA is placed in a cradle for a HotSync process. However, if a user forgets to explicitly close the communications port after using a keyboard that occupied the communications port, the user may be unable to perform an infrared synchronization process. In such cases, the synchronization process cannot be initiated until the situation is resolved. A user unfamiliar with the problem may even have the mistaken impression that the PDA has malfunctioned.\nFurthermore, requiring a user to manually close the communications port also affects power consumption. Battery powered systems, such as PDAs and other hand-held devices, are sensitive to applications that drain excessive power. Thus, PDAs mostly leave their communications ports in a low power standby mode, because an open communications channel may consume significant power. For example, the Palm V personal organizer has a communications port included in the processor that is shared between an infrared communications port and an RS-232 serial communications port. The Palm V hand-held organizer keeps these ports in standby mode to reduce power consumption. A user who forgets to close the communications port after a peripheral device is no longer being used may suffer a significantly decreased battery life, and be forced to recharge or change batteries often.\nEven if a user remembers to close the communications port when the attached peripheral device is no longer being used, the port has to be reopened for further communications with a peripheral device when use is resumed. When the communications port is reopened, the attached peripheral device needs to be identified. Identifying the peripheral device permits the PDA to open appropriate applications and select appropriate communication protocols. Some existing hand-held computers provide limited capability for identifying peripheral devices. For example, in an existing Palm organizer, this is implemented with two pins on a serial port. According to this technique, an interrupt is generated by bringing up one of the pins (e.g., making the pin “high”). An interrupt routine then checks the second pin to determine if it is high or low (e.g., high indicates modem, low indicates cradle). One problem with this technique is that it only allows two different types of devices to wake up the serial port. In other words, this peripheral device identification scheme is limited to identifying one of two devices, for example, a cradle versus a modem. As hand-held computer technology advances and its applications multiply, an increasing number of peripheral devices are becoming available for use with hand-held computers and PDAs, but the existing device identification technique are not capable of distinguishing among the numerous devices.\nSome other approaches to peripheral device identification are capable of identifying more than two types of devices. Typically, such approaches involve performing certain detection activities in a given time period immediately following the initial interrupt generated by the device as described above. According to these approaches, additional input is generated by the peripheral device and received by the hand-held computer during this prescribed time period. For example, the additional input may include one or more specific characters or additional interrupt(s). Based on the additional input, such as by monitoring for the specific inbound character(s) or by counting the additional interrupts, the hand-held computer determines the type of the peripheral device being coupled thereto.\nNevertheless, these identification schemes also have their shortcomings. In particular, these schemes impose a hard limit on the latency within which the initial interrupt generated by the peripheral device must be handled in order for the respective schemes to succeed. Such a time-critical interrupt response requirement is undesirable because the hand-held computer may be engaged in other computations when the initial interrupt is generated and thus may be unable to service the interrupt within the requisite time frame. Under these circumstances, peripheral device identification according to these time-critical schemes would fail.\nThus, in view of the foregoing problems that are inherent in existing peripheral device identification schemes, it would be highly advantageous to provide a peripheral device identification system and method that facilitates extendible identification of various peripheral devices, where the successful operation of the system and method does not require that the initial interrupt generated by the peripheral device be handled within a critical time period.\nAdditionally, a method and system for peripheral device identification should not require complete revamping of existing hand-held computer systems or PDAs. In other words, in implementing a viable method and system for peripheral device identification, components that are well known in the art and are compatible with existing hand-held computer systems need to be used so that the cost of realizing the method and system for peripheral device identification is low. By so doing, the need to incur costly expenditures for retrofitting existing hand-held computer systems or for building custom components is eliminated."}
{"text": "Mercury is a highly toxic element, and globally its discharge into the environment is coming under increasingly strict controls. This is particularly true for power plants and waste incineration facilities. Almost all coal contains small amounts of speciated and elemental mercury along with transition metals (primarily iron) and halogens (primarily chlorine with small amounts of bromine).\nMercury in coal is vaporized in the combustion zone and exits the high temperature region of the boiler entirely as Hg° while the stable forms of halogens are acid gases, namely HCl and HBr. The majority of coal chlorine forms HCl in the flue gas since the formation of elemental or diatomic chlorine is limited due to other dominant flue gas species including water vapor, sulfur dioxide (SO2), nitrogen oxides (NOx) and sulfur trioxide (SO3). By way of example, the Griffin reaction holds that sulfur dioxide, at the boiler temperature range, reacts with elemental or diatomic chlorine to form sulfur trioxide and HCl. Bromine forms both HBr and Br2 at the furnace exit but at temperatures that are important for mercury oxidation, below about 400° C. it is predominantly present in flue gas as Br2. Elemental mercury oxidation occurs primarily via direct halogenation to mercuric chloride and bromide species by both homogeneous gas-phase and heterogeneous surface/gas reactions. For low rank coals with low to medium sulfur and low chlorine and bromine contents, homogeneous gas-phase Hg oxidation reactions are believed to be limited primarily by diatomic Cl2 and Br2 rather than by HCl and HBr due to the slow reaction rate of HCl and HBr. Therefore, though homogeneous gas phase mercury oxidation by diatomic chlorine does occur as the flue gas cools it is not the dominant reaction pathway because insufficient diatomic chlorine is generally present. Rather, heterogeneous reactions controlled by HCl in the cooler regions of the flue gas path past the economizer section and especially occurring within and downstream of the air preheater, on fly ash particles and on duct surfaces are considered to be the primary reaction pathway for oxidation of elemental mercury by chlorine. At cooler flue gas temperatures elemental or diatomic halogens may be formed from HCl and HBr by, for example, a Deacon process reaction. HCl and HBr react with molecular oxygen at cooler flue gas temperatures to form water and diatomic chlorine and bromine, respectively. This reaction is thermodynamically favorable but proceeds only in the presence of metal catalysts that are primarily present on the surface of entrained fly ash particles or on duct surfaces.\nThe U.S. Geological Survey database COALQUAL gives halogen data from analyzed coal specimens. According to this data, U.S. coals have bromine contents between 0 and 160 ppm and the mean and median bromine concentration of the coals are 19 and 12 ppm, respectively, and chlorine contents between 0 and 4,300 ppm and the mean and median chlorine concentration of the coals are 569 and 260 ppm, respectively. Based on the data, lignite and sub-bituminous (e.g., Powder River Basin (“PRB”)) coals are significantly deficient in halogens as compared to average U.S. coals while bituminous coals are higher in halogens than the lower rank coals. For lower rank coals, Hg° is the predominant vapor mercury species.\nVarious methods of augmenting HCl to increase oxidized mercury have been tested at full-scale. Direct addition of halide salts to the coal or injection of halide salts into the boiler has been attempted. There have also been a number of trials of coal blending of low-rank subbituminous coals with higher chlorine bituminous coals. Increased chlorine in the boiler in the form of halide salts or higher chlorine results in an increase of primarily HCl in the flue gas and very limited Cl2. These tests appear to indicate that excess HCl alone does not significantly increase the HgCl2 fraction unless a mechanism exists to make Cl available. Naturally occurring mechanisms that appear to be effective include catalysts in the form of activated carbon or LOI carbon.\nFor lower rank coals, there is thus a need for an effective mercury control methodology."}
{"text": "All cancers are believed to be due to mutations. Testing for myeloid differentiation primary response gene 88 (MYD88) has significant therapeutic and diagnostic value in a range of cancer types, including Waldenström's Macroglobulinemia (WM), diffuse large B-Cell lymphoma (DLBCL), monoclonal gammopathy of unknown significance (MGUS), and splenic marginal zone lymphoma (SMZL), collectively, “MYD88-associated cancers”. MYD88 is an adaptor molecule in a toll-like receptor and interleukin-1 receptor signaling pathway. Mutation in MYD88 results in over-activation of toll-interleukin-receptor pathways, subsequent phosphorylation of IRAK1/4, and release of nuclear factor-kappa-B (NF-κB) drive cell survival and proliferation. It has been demonstrated in DLBCL and WM that inhibition of MYD88 signaling results in decreased NF-κB activity and reduced cell survival. MYD88 mutations are also associated with greater disease burden in patients with DLBCL and poor overall survival following initial and secondary therapy options. MYD88 mutations are detected in 39% of activated B-cell-like (ABC) DLBCL. These mutations are, however, rarely discovered in germinal center B-cell-like (GCB) DLBCL and primary medastinal B-cell lymphoma (PMBL). Therefore, MYD88 mutation status may serve as a surrogate marker for the ABC-subtype. Recent work on ABC-DLBCL and WM has demonstrated increased response to therapy—in both disease types—by combination therapy with a toll-like receptor agonist (IMO-8400) and Rituximab in mouse models. As such, MYD88 mutation status is a useful marker in determining prognosis and in guiding current and future therapy options.\nMYD88 mutations are found in almost all cases of WM and ˜50% of patients with Immunoglobulin M (IgM)-secreting MGUS, while these mutations are rarely detected in patients with SMZL (0-6%) and are absent in multiple myeloma (0%). Differential diagnosis of WM from SMZL and IgM-multiple myeloma is often difficult because of overlapping morphologic, immunophenotypic, cytogenetic, and clinical characteristics. MYD88 is, therefore, a useful marker for accurate diagnosis given its positive mutation status presenting primarily in WM. MYD88's diagnostic, prognostic, and therapeutic power necessitates the development of high-throughput, high-sensitivity assays.\nMost mutations in MYD88 occur at codon L265, converting leucine to proline (L265P); however, mutations at M232, P258, L103, and Q143 have also been reported. Allele-specific (AS) polymerase chain reaction (PCR) based assays have been developed for MYD88 L265P and demonstrated the ability to detect minute fractions of L265P positive cells. However, inherent limitations in this methodology prohibit AS-PCR from detecting variants other than those previously described or specifically designed for this purpose. Thus, the need for an alternative methodology that can be used in routine lab work at larger volumes is desired.\nWild-type blocking PCR (WTB-PCR) using locked nucleic acid (LNA) has demonstrated high sensitivity and versatility in the detection of low percentage mutant populations. By adding an LNA oligo (10-12 NT), complementary to the region of the hotspot, amplification of the WT allele is inhibited, leading to experimentally driven positive selection for mutant alleles. This is accomplished by designing the LNA oligo so that it anneals to the template strand during the primer annealing step of PCR and melts from mutant template DNA—but not WT DNA—during extension. Because a single nucleotide mismatch in the LNA-DNA hybrid greatly decreases its melting temperature, only mutant template DNA is free to complete its extension. Therefore, WT DNA is amplified linearly but mutant DNA is amplified exponentially. Traditional Sanger sequencing can then be performed.\nSanger sequencing has traditionally been the gold standard in testing for both known and unknown somatic mutations. One of the limitations of Sanger sequencing is its limit of detection (˜10-20% mutant allele in a background of WT). This level of sensitivity is inappropriate for detecting low level somatic mutations that may be present in samples from premalignant tissues or patients with few circulating tumor cells, or when bone marrow (BM) is patchy. This also makes assessing residual disease after therapy or detecting emerging resistance mutations during therapy difficult by conventional sequencing alone. By replacing conventional PCR with LNA-mediated wild-type blocking PCR (WTB-PCR) in Sanger sequencing, sensitivities of up to 0.1% mutant allele in a background of WT can be achieved. In WTB-PCR enrichment for mutant alleles is achieved via the addition of a short (˜10-14 NT) LNA oligonucleotide that binds preferentially to WT DNA thereby preventing amplification of WT DNA. The mutant enriched WTB-PCR product can then be sequenced. By blocking WT DNA rather than selecting for specific mutations WTB-PCR allows for enrichment of both known and unknown mutations present in minute cell fractions.\nAmong the most prevalent methods for detecting mutations in small cell fractions are allele-specific PCR (AS-PCR) and real-time quantitative PCR (qPCR). Both are limited by false-positives and the ability to only detect one mutation for which the assay was designed. WTB-PCR, however, allows the user to visualize sequencing traces, which enables the detection of multiple mutation types and can aid in ruling out false positives due to artifacts or deamination events. Next-generation sequencing (NGS) may offer a suitable alternative to conventional sequencing, however, substantially greater costs, complexity, and longer assay time render it an unnecessary option for many disease types with few distinct molecular markers or for monitoring patients on therapy for emerging resistance mutations. Furthermore, high false positive rate when detecting variants with mutant allele frequencies of less than 5% can pose a problem for amplicon-based NGS.\nThe utility of WTB-PCR over AS-PCR lies in its ability to block WT allele amplification rather than amplifying one specific variant allele. In addition, these variant alleles can then be visualized and confirmed via sequencing—contrary to AS-PCR—thereby avoiding false positives. This enrichment of the mutant DNA is particularly useful in testing clinical samples that may contain relatively few neoplastic cells. DLBCL or WM cells may constitute a small percentage of the total cells in bone marrow or peripheral blood samples, leading to false negative results."}
{"text": "Many electronic applications have a need for linear power amplifiers. In the past, power amplifier linearity was often achieved by performing a factory calibration procedure to develop pre-emphasis compensation for the power amplifier to linearize the gain thereof. The compensation information would then be stored within the device carrying the power amplifier for use throughout the life of the amplifier. This factory calibration procedure can be complex and costly to perform for each item in a production line. In addition, this calibration technique only compensates for distortion in the power amplifier under the conditions that exist at the time the calibration is performed and does not provide compensation for effects such as gain variation as a function of junction temperature, gain variation as a function of component aging, non-linearity effects associated with the use of different modulation techniques, and/or other gain variations that might occur during component operation. There is a need for power amplifier linearization techniques that are capable of adapting to changing conditions associated with an amplifier."}
{"text": "A wellhead Christmas tree has a plurality of outlets where valves are placed in connection with these outlets for control of these outlets. These valves are then mounted externally on a large structure and may be subject to damage, for example in the event of falling loads. If the valve is damaged and/or is also completely knocked off, the outlet from the wellhead Christmas tree will be open to the environment. This will result in spillage to the environment and is therefore not desirable. A possible solution to this is to provide an extra barrier valve internally in the actual bore outlet of the wellhead Christmas tree leading out to the valve which may be damaged. This extra barrier valve will have a fail safe closed position. A solution of this kind is described, for example, in the applicant's own Norwegian application NO20080211, entered as PCT application with publication number WO2009102214. This ensures that in the event of damage due to falling loads there is no spillage to the environment.\nIn order to install an extra valve in the bore outlet, the external valve normally has to be uncoupled in order to insert a distance piece between the wellhead Christmas tree and the outer valve, which distance piece is intended for use as a hydraulic chamber for the extra internal valve. Furthermore, the pipe system outside the Christmas tree has to be adapted so as to make room for this distance piece. This distance piece must be included in the assembly as the extra valve which has to be placed in the bore outlet has a control unit/actuator extending slightly beyond the valve and therefore outside the Christmas tree. This control unit/actuator may therefore extend into the main channel in towards the main valve and may possibly hinder closing/opening of the main valve. Thus it is necessary to have a distance piece, and in addition it is also necessary to supply control fluid to the actuator which is to be placed in the bore outlet. This process of installing a distance piece is time-consuming and there is therefore a need for a simpler solution for installation of a valve in the bore outlet in a wellhead Christmas tree."}
{"text": "To meet the demand for wireless data traffic, which has increased since deployment of 4th-generation (4G) communication systems, efforts have been made to develop an improved 5th-generation (5G) or pre-5G communication system. Therefore, the 5G or pre-5G communication system is also called a ‘beyond 4G network’ or a ‘post long-term evolution (LTE) system’.\nIt is considered that the 5G communication system will be implemented in millimeter wave (mmWave) bands, e.g., 60 GHz bands, so as to accomplish higher data rates. To reduce propagation loss of radio waves and increase a transmission distance, a beam forming technique, a massive multiple-input multiple-output (MIMO) technique, a full dimensional MIMO (FD-MIMO) technique, an array antenna technique, an analog beam forming technique, and a large scale antenna technique are discussed in 5G communication systems.\nIn addition, in 5G communication systems, development for system network improvement is under way based on advanced small cells, cloud radio access networks (RANs), ultra-dense networks, a device-to-device (D2D) communication, a wireless backhaul, a moving network, a cooperative communication, coordinated multi-points (CoMP), a reception-end interference cancellation, and the like.\nIn the 5G system, a hybrid frequency shift keying (FSK) and quadrature amplitude modulation (QAM) modulation (FQAM) and a sliding window superposition coding (SWSC) as an advanced coding modulation (ACM) scheme, and a filter bank multi carrier (FBMC) scheme, a non-orthogonal multiple Access (NOMA) scheme, and a sparse code multiple access (SCMA) scheme as an advanced access technology have been developed.\nEnhancement of camera performance and implementation of a personal broadcast service application facilitate live broadcast with a high-definition image of definition equal to or greater than 2 k ultra high definition (UHD) on a client terminal such as a smart phone, and/or the like.\nFor transmitting a high-definition image in real time in a client terminal, a sufficient data rate, e.g., a data rate equal to or greater than 6 Mbps needs to be guaranteed in an uplink channel. However, a bandwidth of a wireless channel which is based on a cellular network (hereinafter, it will be referred to as “cellular channel”) is narrow. Thus, when the number of client terminals in a serving evolved node B (eNB) increases, the data rate per client terminal decreases thereby making it difficult to guarantee a high data rate.\nMeanwhile, a wireless communication system has proposed various technologies for increasing transmission capability, and a typical one is a carrier aggregation (CA) technology. The CA technology has been implemented with various forms, and an inter-heterogeneous eNB CA technology as one of the various forms is a technology in which frequency resources of different communication networks may be aggregated. For example, the inter-heterogeneous eNB CA technology may aggregate frequency resources of an LTE mobile communication system. Further, the inter-heterogeneous eNB CA technology may aggregate a frequency resource of an LTE mobile communication system and a frequency resource of a third generation (3G) mobile communication system, or a frequency resource of an LTE mobile communication system and a frequency resource of a wireless local area network (WLAN), e.g., a Wi-Fi, and/or the like. So, the inter-heterogeneous eNB CA technology may effectively use resources in various communication schemes for increasing a data rate.\nAnother one of the various technologies for increasing the transmission capability proposed in the wireless communication system is a multi-path transmission control protocol (MPTCP). The MPTCP is an example of a new transmission layer protocol which is generated for increasing a data rate using a multi-path. The MPTCP may combine a plurality of transmission paths connecting hosts, e.g., subflows to form one multi-path, e.g., an MPTCP flow and exchange data through the plurality of transmission paths. So, if the MPTCP is used, a preset number of transmission paths may be generated between the hosts.\nFor example, if the MPTCP is used, a data rate sufficient for providing a live video with high definition in real time may be acquired by combining insufficient cellular channel resources, e.g., uplink cellular channel resources. Further, hosts which transmit and receive data based on the MPTCP may distribute traffic to other transmission paths if a congestion situation occurs on a specific transmission path or a connection thereof is released due to various reasons.\nIn a general communication system supporting a multi-path, there is a need for a congestion control scheme for maximizing a data rate between hosts. However, in a current MPTCP standard, a congestion control scheme for an MPTCP is not defined, so a congestion control scheme for a transmission control protocol (TCP) is used. Due to this, congestion control in which a characteristic of an MPTCP is reflected is not performed, so various issues such as fairness violation due to duplicative occupation of a bottleneck in a multi-path, packet reordering due to difference among data rates and difference among round trip time (RTT) of transmission paths, and/or the like may occur.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "Communication devices of the present are capable of as many functions as communications devices from science fiction of the past. As communication moves from desk or wall telephones to cellular phones, voice paging and personal communication systems (better known as PCS), the shrinking size of communications devices coupled with the increase in functionality requires very small, multi-functional controls. As a result, there may be confusion as to how to operate these devices.\nOne area of confusion is signaling; specifically, how to set up a call and how to send further signaling tones through the network once the call is set up. Users are used to one paradigm for a dual-tone, multi-frequency (DTMF) or A touch-tones signaling (referred to herein as DTMF mode or DTMF signaling). Most anyone can use a DTMF keypad to generate signaling that performs tasks as diverse as dialing a telephone number to registering for classes and banking. DTMF signaling presumes that a pair of tones is generated for each key or button pushed and, once the key is pushed, the signal it represents cannot be changed. Call routing menu systems are examples of systems that expect one button, immediate generation of DTMF signaling in this paradigm.\nWireless telephones (also called cell phones, mobile stations and mobile telephones) use a different signaling paradigm. In this paradigm, the user enters digits, which are usually displayed on a screen, and then the user presses a \"send\" or \"enter\" key to send the digits into the network as a message. Commonly, there is a \"backspace\" or \"erase\" key that permits the user to change one or more entered digits, and a \"clear\" key to remove all entered digits before the signaling message is sent. This form of signaling will be referred to herein as \"dialed digit\" or \"DD\" mode. Once the signaling message is sent, some wireless telephones then cue to a DTMF mode, wherein the user may send DTMF signals through the telecommunication network to perform tasks such as banking, etc.\nIn the dual mode telephones that have two paradigms operating at any given time, however, it is not always clear to the user operating the telephone which mode a communication device is in (DTMF or DD), and the user cannot control the mode. Furthermore, there are times in the DTMF mode when it is desirable to be able to correct digits before they are sent into the network (i.e., when entering a string of digits representing a credit card number). Consequently, users of new communication devices can easily be confused by the two different signaling paradigms."}
{"text": "1. Field of the Invention\nThe present invention relates to a heating resistor element component (thermal head) which is used in a thermal printer typically mounted onto a compact information equipment terminal such as a compact handy terminal, and is used for performing printing on a thermal recording medium through selective driving of a plurality of heating elements based on print data.\n2. Description of the Related Art\nRecently, thermal printers have been widely used in compact information equipment terminals. The compact information equipment terminals are driven by a battery, which leads to strong demands for electric power saving of the thermal printers. Accordingly, there have been growing demands for heating resistor element components having high heating efficiency.\nAs to increasing efficiency of the thermal head, there is known a method of forming a hollow portion in a lower layer of a heating resistor (for example, see JP 2007-83532 A). Among an amount of heat generated in the heating resistor, an amount of upper-transferred heat which is transferred to a wear-resistant layer formed above the heating resistor becomes larger than an amount of lower-transferred heat which is transferred to a supporting substrate located under the heating resistor, and thus energy efficiency required during the printing can be sufficiently obtained.\nHowever, in the heating resistor element component disclosed in JP 2007-83532 A, a heat accumulating layer (insulating film) provided on a surface of the supporting substrate is required to have a thickness to some extent in terms of a mechanical strength, which imposes a limitation on reducing a thickness (size in a height direction) of the heating resistor element component."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to a coating that is applied to a surface. More particularly, embodiments of the present invention relate to a protective coating that is applied to the underwater portion of a marine vessel so as to inhibit the growth of marine foulants.\n2. Description of the Related Art\nMarine vessels that reside in a water environment over certain lengths of time can accumulate biological growth, known as foulants, on those surfaces that are in contact with the water. Diverse species of hard and soft fouling organisms, such as barnacles, zebra mussels, algae, and slime, form colonies on the underwater surfaces of the vessel, particularly when a vessel is docked, because each requires a permanent anchorage in order to mature and reproduce. Marine growth fouling adds weight to a ship, increases the amount of fuel consumed, and reduces its speed.\nHistorically, to combat the growth of marine foulants, the underwater surfaces of ships have been coated with antifoulant paints, which often include toxic materials to inhibit biological growth. The antifoulant paints may degrade and break down over time, releasing the toxic materials from the marine vessel into the surrounding water. These toxic materials may include volatile organic compounds (VOCs) and hazardous air pollutants (HAPs). The International Maritime Organization and the United States Environmental Protection Agency have enacted regulations and standards that restrict the emission of VOCs and HAPs from antifoulant paints. The decomposition and break down of the antifoulant paint results in reduced efficacy of the protection afforded by the antifoulants, thereby requiring reapplication of the paint in a relatively short time. Thus, a coating material is required that can be applied to the underwater surfaces of a marine vessel which repels the growth of fouling organisms on such surfaces and has an extended lifetime without releasing significant amounts of toxic materials into the environment."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical coupling device and a master used in fabricating the optical coupling device, more particular, to an integrated type optical coupling device and a master used in fabricating the optical coupling device that an optical alignment between a narrow-pitch multi-channel optical waveguide and an optical fiber array can be easily accurately performed and the fabricating cost thereof can be reduced.\n2. Description of the Prior Art\nRecently, as an information communication industry is developed, data transmission using a communication network such as Internet is actively increased, thereby large scale, high speed, and high density of a transmission/exchange system for increasing the transmitting amount and the transmitting speed of data has been demanded. According to the demand, the research for integrating the channel to 40 channels greater than the existing 8 or 16 channels in a same area has been progressed. At this case, since the pitch between the channels is reduced from 250 μm to 125 μm, it is difficult to fabricate the system. Accordingly, the implement of the optical device which can easily integrate the channel, can increase optical coupling efficiency with the optical fiber, and can reduce the fabricating cost thereof is needed.\nGenerally, the optical coupling device for optically coupling with the optical fiber and electrically or optically controlling an optical signal comprises modules formed with a plurality of V-shaped grooves for mounting an optical fiber array, a waveguide element coupled between the modules and formed with a multi-channel optical waveguide, and a controlling means for electrically or optically controlling the optical signal transmitted through the optical waveguide.\nIn the conventional optical coupling device having the above-mentioned structure, the module and the waveguide element are separately fabricated and are coupled with each other. Accordingly, first, optical coupling loss can be generated due to the variation if a thermal or mechanical impact is applied thereto. Second, since the dynamic stability of the optical fiber is not secured due to the space of the lower portion when the optical fiber is aligned in the V-shaped groove, the length of the V-shaped groove must be relatively long, thereby it is difficult to reduce the size of the structure. Also, third, an expensive optical fiber aligning device must be used in order to make a precise optical alignment between the optical waveguide and the optical fiber array. Fourth, since the optical fiber is fixed by an adhesive to make the active optical alignment with the waveguide element after the intensity of the transmitted light is checked by the channel, the time and technical cost required for aligning the optical fiber array is proportionally increased as the number of the channels is increased.\nOn the other hand, a master (a metallic pattern) is used in manufacturing the module mounted with the optical fiber and the waveguide element. The multi-step master is fabricated by a patterning process using several masks each having a different exposed location. Accordingly, there is a merit that multi-step structure having a different height proportional to the exposed number is made, but the masks as same number as the exposed number must be fabricated, an accurate location alignment of the mask is need in each of the exposing steps, and it is difficult to accurately adjust the time for exposing the photoresist film and the height of the fine structure.\nAs an alternative method, a method for fabricating a master, applying the photoresist film thereon, and performing the plating to implement a fine structure having various heights on the same surface is suggested. In this method, since the photoresist film is formed on the master and then is exposed by using the mask, a precise mask aligning technique is needed, the mask fabricating processes as same number as the number of the steps of the mask must be repeatedly performed. Above all, since a precise dimension control determines the precision of the master which is finally fabricated in the master fabricating process, a precise machining technique is needed. In addition, there is a problem that the fine structure fabricated previously is damaged when forming the photoresist on the fabricated master."}
{"text": "1. Field of the Invention\nThe present invention relates to a gas compressor which is to be favorably assembled into a refrigerating machine or an air conditioner as a refrigerant compressor. Particularly, the invention favorably relates to a motor-driven gas compressor into which an electric motor is assembled as a driving source for a gas-compressing mechanism.\n2. Related Art Statement\nHeretofore, a vane rotary type gas compressor is used as the gas compressor of this type (For example, see JP-A 2003-254244). As shown in FIG. 4, a compressing mechanism 3 is housed in a housing 2 of such a vane rotary type gas compressor 1. The housing 2 is provided with a suction port and a discharge port not shown. The compressing mechanism 3 has a cylinder 4 and a rotor 5 rotatably housed in the cylinder. The cylinder 4 defines a cylinder chamber 6 having an elliptical cross-sectional shape. In the cylinder chamber 6, the rotor 5 is provided with vanes 7 which divide the cylinder chamber 6 into a plurality of compressing chambers 6a in a circumferential direction. The volume of each of the compressing chambers 6a increases or decreases with the rotation of the rotor 5.\nWhen the volume of the compressing chamber 6a increases with the rotation of the rotor 5, a refrigerant is sucked into a suction chamber (not shown) inside the housing 2 through the suction port, and further is sucked from this suction chamber into the compression chamber 6a via a suction passage (not shown) extending in the cylinder 4 and a suction inlet 8 opened to the cylinder chamber 6. On the other hand, when the volume of the compressing chamber 6a decreases with the rotation of the rotor 5, the refrigerant inside the compression chamber is compressed. The compressed refrigerant is guided, via the discharge outlet 9 opened to the cylinder chamber 6, to a discharge space 9a formed by a cut portion defined between the cylinder 4 and the housing surrounding the cylinder 2, into the discharge port communicating with the discharge space. In the discharge space 9a is provided a check valve mechanism 9b for preventing the compressed refrigerant from flowing back into the compressing chamber 6a. \nAlthough not shown, an electric motor is arranged as a driving force for the rotor inside the suction chamber in communication with the suction port of the housing so that the motor may be cooled with the refrigerant flowing in the housing. The actuation of the electric motor operates the gas-compressing mechanism, the refrigerant is sucked into the gas compressor via the suction port of the housing from a compressor (not shown) as mentioned above, and the refrigerant compressed by the gas compressor is fed to a condenser through the discharge port formed in the housing under pressure.\nIn the conventional gas compressor, the suction inlet 8 is opened at an end wall of the cylinder chamber 6 for taking the refrigerant into the cylinder chamber 6, and the refrigerant is sucked from the suction passage into the compressing chamber 6a through the suction inlet 8 opened at the end face of the cylinder chamber 6. On the other hand, the refrigerant compressed in the compressing chamber 6a is discharged into the discharge space 9 formed between the cylinder 4 and the housing 2 covering the cylinder, through the discharge outlet 9a opened in the peripheral wall of the cylinder chamber 6. Therefore, the high pressure of the compressed refrigerant acts as an internal pressure upon a portion of the housing 2 covering the outer peripheral portion of the cylinder of the housing. Consequently, in order to withstand such a high pressure, the housing 2 is made thick, so that the increased thickness may cause the larger dimension and increased weight of the housing.\nFurther, according to the conventional structure in which the discharge pressure acts upon the outer peripheral portion of the cylinder, it was not easy to assuredly prevent internal leakage from occurring due to the leakage of the high pressure at the outer peripheral portion of the cylinder toward a low-pressure side such as the suction passage."}
{"text": "Conventional cameras transmit the total field of view within one channel and are limited in terms of their miniaturization. In mobile devices such as smart phones, for example, two cameras are employed which are oriented in and counter to the direction of the surface normal of the display.\nTherefore, what would be desirable is a concept which enables miniaturized devices to capture a total field of view while ensuring high image quality."}
{"text": "1. Field of the Invention\nThe present invention pertains to the art of refrigerators and, more particularly, to an air flow arrangement for a refrigerator assembly used to store food items in a controlled humidity and temperature environment.\n2. Discussion of the Prior Art\nIn the art of refrigerators, particularly household refrigerators, it is often desirable to create varying humidity and/or temperature storage zones to enhance the preservation of different food items. For instance, it is common to accommodate the storage requirements for certain food items, such as dairy products, meats, fruits and vegetables, by forming separately enclosed storage areas within a fresh food compartment. In most instances, these storage areas are designed to be maintained at temperatures which are different from the temperature of the remainder of the fresh food compartment.\nIn at least the case of fruits and vegetables, it is typically desirable to isolate these food items from direct contact with a flow of cooling air, especially any cold air flowing into the fresh food compartment from a freezer compartment of the refrigerator, mainly because this cold air can be fairly dry. Therefore, in order to isolate the fruits and vegetables from the desiccating effects of the cold air so as to maintain the moisture content of the fruits and vegetables, it has heretofore been proposed to provide a specialized storage receptacle, such as a crisper, within a refrigerator fresh food compartment. A crisper generally takes the form of a slidable bin which is sealed to maintain a relatively high humidity level, while the walls of the bin are chilled to establish a desirable temperature within the bin.\nMany different food storage compartment designs have been proposed in the art in an attempt to establish and maintain effective humidity and temperature conditions within the compartment while attempting to avoid the development of condensation. However, there still exists a need for an improved storage compartment system which can maintain a desired humidity level, accurately control the temperature and minimize the tendency for condensation within the compartment. In designing a specialty food storage compartment for a refrigerator, consideration must be taken of the manner in which a flow of cooling air is effectively and efficiently supplied to and exhausted from the food, vegetable, meat or the like bin."}
{"text": "1. Field of the Invention\nThe present invention relates to a zoom lens device whose lens barrel elements rotate relative to each other for changing the focal length, i.e. for zooming. The present invention relates more particularly to a zoom lens device having a zooming position detector mounted in a compact fashion.\n2. Background Arts\nA zoom lens device has at least front and rear lens groups whose positions are changed relative to a film surface and relative to each other in the direction of the optical axis of the zoom lens device. Since there is a clearance between barrels that hold the front and rear lens groups respectively, extraneous light would enter inside the zoom lens device through the clearance without any light-shielding device. JPY 5-31634 discloses a light-shielding device for a zoom lens device, which is constituted of a light-shielding ring having a resilient lip around its inner rim. The light-shielding ring is placed behind a face flange surrounding an opening of a front wall of a camera body such that the resilient lip stays in contact with the outer periphery of a lens barrel or lens cover frame that moves with the front lens group in the axial direction through the opening of the front wall.\nThe light-shielding ring disclosed in JPY 5-31634 is useful for a clearance around a lens barrel that moves in the axial direction relative to a fixed barrel. However, the light-shielding ring is not preferable for a lens device where a front lens frame is rotatably held in a barrel, and a front cover member having an opening for exposing the front lens is attached to a front end of the barrel, and a clearance is provided between the front lens frame and the front cover member for allowing the front lens frame to rotate relative to the front cover member. If the light-shielding ring is used for the clearance between the front cover member and the front lens frame, the outer periphery of the front lens frame would rub against the lip of the light-shielding ring as the front lens frame rotates relative to the front cover member. The friction between the front lens frame and the front cover member would raise the necessary driving force for the lens device. The lip of the light-shielding ring would sooner be worn out or heat-deformed by the frequent friction. Besides, the light-shielding ring increases the requisite number of parts necessary for the lens device.\nOn the other hand, it is known in the art to provide a zooming position detector in a zoom lens device. For example, JPA 50-36118 discloses a zooming position detector consisting of an encoder plate and a brush device. The encoder plate is coupled to a lens barrel of a zoom lens device through gears and a lead screw, such that the encoder plate moves in the axial direction of the zoom lens device as the lens barrel rotates for zooming. As the encoder plate moves in the axial direction, the brush device brushes conductor patterns on the encoder plate detecting signals corresponding to the rotational angle of the lens barrel. Japanese Utility Model Registration No. 2521469 discloses a zooming position detector wherein a code plate having conductor patterns are tightly provided on an outer periphery of a cam ring that rotates for zooming, whereas a brush device brushing the conductor patterns is secured to a stationary frame outside the cam ring.\nEither of the prior zooming position detectors needs a mounting space outside the movable lens barrel, which inevitably enlarges the whole size of the camera."}
{"text": "The present invention relates to a magnetic resonance signal acquiring method and apparatus, a recording medium and a magnetic resonance imaging apparatus, and more particularly to a magnetic resonance signal acquiring method and apparatus for effecting RF (radio frequency) excitation on spins within an object and acquiring a navigator echo and an imaging echo, a recording medium recorded with a program for causing a computer to perform such a magnetic resonance signal acquiring function, and a magnetic resonance imaging apparatus comprising such magnetic resonance signal acquiring means.\nIn a magnetic resonance imaging (MRI) apparatus, an object to be imaged is carried into an internal space of a magnet system, i.e., a space in which a static magnetic field is generated; a gradient magnetic field and a high frequency magnetic field are applied to generate magnetic resonance signals from spins within the object; and a tomographic image is constructed based on the received signals. The magnetic resonance signal for use in constructing the tomographic image is sometimes referred to as an imaging echo.\nSince the spin frequency is proportional to the magnetic field strength, the spin frequency varies as the static magnetic field strength varies. If the temporal variation in the spin frequency is considerably large, the current value of the spin frequency is determined before effecting RF excitation for generating an imaging echo, and the RF excitation for generating the imaging echo is effected by the high frequency magnetic field having a frequency equal to the current value of the spin frequency.\nIn order to evaluate the current value of the spin frequency, the spins are RF-excited separately to measure an FID (free induction decay) signal. The FID is sometimes referred to as navigator echo. The spin frequency is determined from the navigator echo by a calculation.\nFIG. 1 is a time chart of navigator echo collection, a calculation on the navigator echo and imaging echo collection. The timing of gradient magnetic field application is omitted in the drawing.\nAs shown, the navigator echo collection and the imaging echo collection are separately performed during two consecutive periods. The length of each period is TR (repetition time).\nIn the earlier period, RF excitation for generating a navigator echo is effected at a time tl. Next, from a time t2 to a time t3, navigator echo collection is performed.\nNext, from a time t3 to a time t5, a calculation is performed based on the collected echo data. The spin frequency is determined by the calculation. The result of the calculation is reflected on the frequency of the RF excitation in the later period.\nIn the later period, RF excitation for generating an imaging echo is effected at a time t6. The frequency of the RF excitation is made equal to the spin frequency determined by the calculation in the earlier period. Next, from a time t7 to a time t8, imaging echo collection is performed.\nBy repeating the above operation, imaging echoes for a plurality of views are sequentially collected. Each of the imaging echoes is given different phase encoding for each view by the gradient magnetic field, which is omitted in the drawing.\nIn the above operation, the calculation on the navigator echo should be completed within an interval after the navigator echo collection up to the beginning of the next period. However, this interval is short because it is the remainder of the period TR after subtraction of the time required for the RF excitation and navigator echo collection, and the time available for performing the calculation is insufficient. Accordingly, the calculation must be sped up to complete the calculation within the interval, or, if this is not possible, the period TR must be extended.\nTherefore, it is an object of the present invention to provide a magnetic resonance signal acquiring method and apparatus in which the time available for a calculation on a navigator echo is large, and a recording medium recorded with a program for causing a computer to perform such a magnetic resonance signal acquiring function, and a magnetic resonance imaging apparatus comprising such magnetic resonance signal acquiring means.\n(1) The present invention, in accordance with one aspect for solving the aforementioned problem, is a magnetic resonance signal acquiring method that executes the steps of effecting RF excitation on spins within an object and acquiring a navigator echo; performing a calculation on said acquired navigator echo and adjusting the frequency of RF excitation based on the result of the calculation; and effecting RF excitation on the spins within the object and acquiring an imaging echo, which method is characterized in comprising: making the sampling frequency for acquiring said navigator echo higher than the sampling frequency for acquiring said imaging echo.\n(2) The present invention, in accordance with another aspect for solving the aforementioned problem, is a magnetic resonance signal acquiring apparatus for acquiring magnetic resonance signals from an object using a static magnetic field, a gradient magnetic field and a high frequency magnetic field, characterized in that the apparatus comprises: echo acquiring means for effecting RF excitation on spins within the object and acquiring a navigator echo, and for effecting RF excitation on the spins within the object and acquiring an imaging echo; calculating means for performing a calculation on said acquired navigator echo; frequency adjusting means for adjusting the frequency of RF excitation based on the result of said calculation; and control means for making the sampling frequency for said echo acquiring means to acquire the navigator echo higher than the sampling frequency for said echo acquiring means to acquire the imaging echo.\n(3) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a recording medium that recorded in a computer-readable manner with a program for causing a computer to execute the functions of effecting RF excitation on spins within an object and acquiring a navigator echo; performing a calculation on said acquired navigator echo and adjusting the frequency of RF excitation based on the result of the calculation; and effecting RF excitation on the spins within the object and acquiring an imaging echo, which recording medium is characterized in that the program causes the computer to: make the sampling frequency for acquiring said navigator echo higher than the sampling frequency for acquiring said imaging echo.\n(4) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance imaging apparatus for acquiring magnetic resonance signals from an object using a static magnetic field, a gradient magnetic field and a high frequency magnetic field, and constructing an image based on said acquired magnetic resonance signals, characterized in that means for acquiring said magnetic resonance signals comprises: echo acquiring means for effecting RF excitation on spins within the object and acquiring a navigator echo, and for effecting RF excitation on the spins within the object and acquiring an imaging echo; calculating means for performing a calculation on said acquired navigator echo; frequency adjusting means for adjusting the frequency of RF excitation based on the result of said calculation; and control means for making the sampling frequency for said echo acquiring means to acquire the navigator echo higher than the sampling frequency for said echo acquiring means to acquire the imaging echo.\nAccording to the invention in these aspects, since the sampling frequency for acquiring the navigator echo is made higher than the sampling frequency for acquiring the imaging echo, the time required for the navigator echo data acquisition can be reduced relative to the time required for the imaging echo data acquisition. Thus, the time available for the calculation on the navigator echo can be increased.\n(5) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance signal acquiring method that executes the steps of effecting RF excitation on spins within an object and acquiring a navigator echo; performing a calculation on said acquired navigator echo and adjusting the frequency of RF excitation based on the result of the calculation; and effecting RF excitation on the spins within the object and acquiring an imaging echo, which method is characterized in comprising: making the number of sampling points for acquiring said navigator echo fewer than the number of sampling points for acquiring said imaging echo.\n(6) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance signal acquiring apparatus for acquiring magnetic resonance signals from an object using a static magnetic field, a gradient magnetic field and a high frequency magnetic field, characterized in that the apparatus comprises: echo acquiring means for effecting RF excitation on spins within the object and acquiring a navigator echo, and for effecting RF excitation on the spins within the object and acquiring an imaging echo; calculating means for performing a calculation on said acquired navigator echo; frequency adjusting means for adjusting the frequency of RF excitation based on the result of said calculation; and control means for making the number of sampling points for said echo acquiring means to acquire the navigator echo fewer than the number of sampling points for said echo acquiring means to acquire the imaging echo.\n(7) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a recording medium recorded in a computer-readable manner with a program for causing a computer to execute the functions of effecting RF excitation on spins within an object and acquiring a navigator echo; performing a calculation on said acquired navigator echo and adjusting the frequency of RF excitation based on the result of the calculation; and effecting RF excitation on the spins within the object and acquiring an imaging echo, which recording medium is characterized in that the causes the computer to: make the number of sampling points for acquiring said navigator echo fewer than the number of sampling points for acquiring said imaging echo.\n(8) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance imaging apparatus for acquiring magnetic resonance signals from an object using a static magnetic field, a gradient magnetic field and a high frequency magnetic field, and constructing an image based on said acquired magnetic resonance signals, characterized in that means for acquiring said magnetic resonance signals comprises: echo acquiring means for effecting RF excitation on spins within the object and acquiring a navigator echo, and for effecting RF excitation on the spins within the object and acquiring an imaging echo; calculating means for performing a calculation on said acquired navigator echo; frequency adjusting means for adjusting the frequency of RF excitation based on the result of said calculation; and control means for making the number of sampling points for said echo acquiring means to acquire the navigator echo fewer than the number of sampling points for said echo acquiring means to acquire the imaging echo.\nAccording to the invention in these aspects, since the number of sampling points for acquiring the navigator echo is made fewer than the number of sampling points for acquiring the imaging echo, the time required for the navigator echo data acquisition can be reduced relative to the time required for the imaging echo data acquisition. Thus, the time available for the calculation on the navigator echo can be increased.\n(9) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance signal acquiring method that executes the steps of effecting RF excitation on spins within an object and acquiring a navigator echo; performing a calculation on said acquired navigator echo and adjusting the frequency of RF excitation based on the result of the calculation; and effecting RF excitation on the spins within the object and acquiring an imaging echo, which method is characterized in comprising: making the sampling frequency for acquiring said navigator echo higher than the sampling frequency for acquiring said imaging echo; and making the number of sampling points for acquiring said navigator echo fewer than the number of sampling points for acquiring said imaging echo.\n(10) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance signal acquiring apparatus for acquiring magnetic resonance signals from an object using a static magnetic field, a gradient magnetic field and a high frequency magnetic field, characterized in that the apparatus comprises: echo acquiring means for effecting RF excitation on spins within the object and acquiring a navigator echo, and for effecting RF excitation on the spins within the object and acquiring an imaging echo; calculating means for performing a calculation on said acquired navigator echo; frequency adjusting means for adjusting the frequency of RF excitation based on the result of said calculation; and control means for making the sampling frequency for said echo acquiring means to acquire the navigator echo higher than the sampling frequency for said echo acquiring means to acquire the imaging echo, and for making the number of sampling points for said echo acquiring means to acquire the navigator echo fewer than the number of sampling points for said echo acquiring means to acquire the imaging echo.\n(11) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a recording medium recorded in a computer-readable manner with a program for causing a computer to execute the functions of effecting RF excitation on spins within an object and acquiring a navigator echo; performing a calculation on said acquired navigator echo and adjusting the frequency of RF excitation based on the result of the calculation; and effecting RF excitation on the spins within the object and acquiring an imaging echo, which recording medium is characterized in that the program causes the computer to: make the sampling frequency for acquiring said navigator echo higher than the sampling frequency for acquiring said imaging echo; and make the number of sampling points for acquiring said navigator echo fewer than the number of sampling points for acquiring said imaging echo.\n(12) The present invention, in accordance with still another aspect for solving the aforementioned problem, is a magnetic resonance imaging apparatus for acquiring magnetic resonance signals from an object using a static magnetic field, a gradient magnetic field and a high frequency magnetic field, and constructing an image based on said acquired magnetic resonance signals, characterized in that means for acquiring said magnetic resonance signals comprises: echo acquiring means for effecting RF excitation on spins within the object and acquiring a navigator echo, and for effecting RF excitation on the spins within the object and acquiring an imaging echo; calculating means for performing a calculation on said acquired navigator echo; frequency adjusting means for adjusting the frequency of RF excitation based on the result of said calculation; and control means for making the sampling frequency for said echo acquiring means to acquire the navigator echo higher than the sampling frequency for said echo acquiring means to acquire the imaging echo, and for making the number of sampling points for said echo acquiring means to acquire the navigator echo fewer than the number of sampling points for said echo acquiring means to acquire the imaging echo.\nAccording to the invention in these aspects, since the sampling frequency for acquiring the navigator echo is made higher than the sampling frequency for acquiring the imaging echo, and the number of sampling points for acquiring the navigator echo is made fewer than the number of sampling points for acquiring the imaging echo, the time required for the navigator echo data acquisition can be reduced relative to the time required for the imaging echo data acquisition. Thus, the time available for the calculation on the navigator echo can be increased.\nTherefore, the present invention can provide a magnetic resonance signal acquiring method and apparatus in which the time available for a calculation on a navigator echo is long, and a recording medium on which a program is recorded for causing a computer to perform such a magnetic resonance signal acquiring function, and a magnetic resonance imaging apparatus comprising such magnetic resonance signal acquiring means."}
{"text": "Limitations in modulating bandwidth and efficiency in radio frequency (RF) power amplifiers prevent reductions in size, weight, power and cooling. Class-D power amplifiers offer the highest theoretical efficiency and bandwidth of known power amplifiers. The necessity of driving a floating, high side switch has hindered their widespread implementation.\nSwitching power converters utilizing power switching devices in a half-bridge or full-bridge configuration are commonly operated in conditions where the active devices pass reverse current. In many cases, a schottky diode is placed in parallel with each active device to permit conduction of reverse current in a path other than the main switching devices. These schottky diodes add parasitic capacitance and inductance.\nConsequently, it would be advantageous if an apparatus existed that is suitable for power switching without complicated drive circuitry."}
{"text": "The present invention relates generally to Java programming environments, and more particularly, to techniques suitable for determining hashcode values associated with Java objects.\nOne of the goals of high level languages is to provide a portable programming environment such that the computer programs may easily be ported to another computer platform. High level languages such as “C” provide a level of abstraction from the underlying computer architecture and their success is well evidenced from the fact that most computer applications are now written in a high level language.\nPortability has been taken to new heights with the advent of the World Wide Web (“the Web”) which is an interface protocol for the Internet which allows communication between diverse computer platforms through a graphical interface. Computers communicating over the Web are able to download and execute small applications called applets. Given that applets may be executed on a diverse assortment of computer platforms, the applets are typically executed by a Java virtual machine.\nRecently, the Java programming environment has become quite popular. The Java programming language is a language that is designed to be portable enough to be executed on a wide range of computers ranging from small devices (e.g., pagers, cell phones and smart cards) up to supercomputers. Computer programs written in the Java programming language (and other languages) may be compiled into Java Bytecode instructions that are suitable for execution by a Java virtual machine implementation. The Java virtual machine is commonly implemented in software by means of an interpreter for the Java virtual machine instruction set but, in general, may be software, hardware, or both. A particular Java virtual machine implementation and corresponding support libraries together constitute a Java runtime environment.\nComputer programs in the Java programming language are arranged in one or more classes or interfaces (referred to herein jointly as classes or class files). Such programs are generally platform, i.e., hardware and operating system, independent. As such, these computer programs may be executed without modification on any computer that is able to run an implementation of the Java runtime environment.\nObject-oriented classes written in the Java programming language are compiled to a particular binary format called the “class file format.” The class file includes various components associated with a single class. These components can be, for example, methods and/or interfaces associated with the class. In addition, the class file format can include a significant amount of ancillary information that is associated with the class. The class file format (as well as the general operation of the Java virtual machine) is described in some detail in The Java Virtual Machine Specification, Second Edition, by Tim Lindholm and Frank Yellin, which is hereby incorporated herein by reference.\nFIG. 1A shows a progression of a simple piece of a Java source code 101 through execution by an interpreter, the Java virtual machine. The Java source code 101 includes the classic Hello World program written in Java. The source code is then input into a Bytecode compiler 103 that compiles the source code into Bytecodes. The Bytecodes are virtual machine instructions as they will be executed by a software emulated computer. Typically, virtual machine instructions are generic (i.e., not designed for any specific microprocessor or computer architecture) but this is not required. The Bytecode compiler outputs a Java class file 105 that includes the Bytecodes for the Java program. The Java class file is input into a Java virtual machine 107. The Java virtual machine is an interpreter that decodes and executes the Bytecodes in the Java class file. The Java virtual machine is an interpreter, but is commonly referred to as a virtual machine as it emulates a microprocessor or computer architecture in software (e.g., the microprocessor or computer architecture may not exist in hardware).\nFIG. 1B illustrates a simplified class file 100. As shown in FIG. 1B, the class file 100 includes a constant pool 102 portion, interfaces portion 104, fields portion 106, methods portion 108, and attributes portion 110. The methods portion 108 can include or have references to several Java methods associated with the Java class which is represented in the class file 100.\nAs is known in the art, Java objects can be uniquely identified by a hashcode value. During the execution of Java programs, there is often a need to determine hashcode values for Java objects. Conventionally, the Java method “/java/lang/object.hashcode” is invoked by the virtual machine to determine these hashcode values. One problem with this approach is that there is an overhead associated with the invocation of a Java method. In other words, Java method invocation requires several operations to be performed. These operations include: locating the appropriate method to be invoked, creating a frame to be placed on the execution stack, and restoring the previous frame on the stack.\nMoreover, the cost associated with determining hashcode values is quite high, since the Java method has to be invoked time and time again. In other words, the operations needed to invoke a method have to be performed several times during the execution of a Java program. This, of course, can result in a grossly inefficient use of system resources. In some circumstances, particularly in systems with limited computing power and/or memory, this inefficient use of resources is a serious disadvantage.\nIn view of the foregoing, there is a need for improved techniques to determine Java hashcode values which are used to identify Java objects."}
{"text": "Resistance of tumor cells to chemotherapeutic drugs is an important problem in the clinical management of cancer. One type of resistance is characterized by the cross resistance to a wide variety of chemotherapeutic agents with no major structural similarities or similar modes of action. This phenomenon is termed multiple drug resistance. The evidence suggests that this form of resistance is due to the presence in tumor cells of a 170 kD trans-membrane protein that is termed P-glycoprotein. It is thought that resistance is conferred by the ability of P-glycoprotein to actively transport chemotherapeutic agents out of the cells and thereby lower the intracellular concentrations of such drugs to non-toxic levels. Clinically this translates into a reduction of the therapeutic index of such drugs to an ineffective level. It has been shown that MDR1 mRNA levels are elevated in untreated, intrinsically drug resistant tumors including those derived from the colon, kidney, adrenal gland, liver and pancreas as well as some cancers at relapse after chemotherapy including breast cancer, neuroblastoma, ALL and nodular poorly differentiated lymphoma.\nChemotherapeutic agents in clinical use to which multi-drug resistance has been observed include doxorubicin, vinblastine, vincristine, taxol, duanomycin, etopiside, teniposide and actinomycin D.\nIn the early 1980's it was discovered that multiple drug resistance tumor cells could be resensitized to chemotherapeutic agents with the anti-hypertensive, verapamilo Since this time verapamil, and several other agents including R-verapamil and trifluoroperizine have been used as resensitizing agents for the treatment of cancer in the clinic. Their activity has been limited by toxic side effects.\nThe present invention describes various bicyclic amines that are effective in resensitizing multiple drug resistant cells to chemotherapeutic agents such as doxorubicin, vincristine and bisantrene. This resensitization has the effect of reducing the amount of chemotherapeutic agent neccessary to kill those cells. In addition this potentiation of activity of the chemotherapeutic agent with the compounds of this invention makes it possible to treat multi-drug resistant tumors in warm blooded animals.\nAn object of this invention is provide a method for increasing the sensitivity of multi-drug resistant cells by the use of compounds in this invention.\nAn additional object of this invention is to provide a method for treating tumors that are either intrinsically or extrinsically multi-drug resistant in warm blooded animals by the administration of a therapeutically effective dose of a compound of this invention prior to, concurrent or after the administration of a therapeutically effective dose of an antitumor chemotherapeutic agent."}
{"text": "In particular, yet not restrictive manner, the present application concerns a device and a method of producing a laminate floor panel with a curve shaped and compressed edge. However, the application is as well applicable to building panels in general. More particularly, the application relates to the type of floor panels with a wood fibre based core, e.g. HDF, MDF, particleboard, plywood, a decorative layer and a mechanically locking system for horizontal and/or vertical locking produced by the method and device.\nA device, method and floor panel of this type is presented in WO2006/088417, which discloses a device and method to compress the edges of a laminate floor panel to provide a laminate floor panel with curve shaped edges looking like a solid wood floor panel with bevelled edges.\nThe contents of WO2006/088417 and U.S. Ser. No. 10/906,356 are hereby incorporated herein by reference."}
{"text": "A hydrostatic actuator of this type is known from German Patent No. 197 00 935 A1. In accordance with this document, a piston of a master cylinder is driven by an electric motor and the rotary movement of the electric motor is converted into an axial movement by a gearing mechanism in the form of a worm drive. The piston is driven by an eccentric pin arranged on the worm gear.\nIn addition, European Patent No. 0 320 621 A1 discloses a planetary rolling-contact gear system that includes a threaded spindle, a sleeve arranged coaxially with the threaded spindle, and planetary gears rolling off between the two. The planetary gearing mechanism is driven to rotate by an electric motor and converts a fast rotary movement into a slower axial movement.\nThus, there is a long-felt need for a hydrostatic actuator that is preferably used in a motor vehicle and has an increased efficiency, reduced required installation space, and increased operating pressure of the master cylinder."}
{"text": "Driver or buffer circuits are often used to help isolate outputs and match output impedance to expected load impedances. An example of one known differential output driver 100 is illustrated in FIG. 1. The driver 100 includes a pair of transistors Q1, Q2 in a push-pull configuration driving a resistive load RL. A direct current (DC) bias circuit 102 maintains a constant DC current through the transistors Q1, Q2. The DC bias circuit includes a pair of transistors M1, M2 with gates tied to the gate of a transistor M0 driven by a current source I0. Unfortunately, this type of driver can generate undesirable amounts of distortion.\nAs the input voltages Vp and Vm change during operation of the driver 100, the output voltages of Q1 and Q2, Vop and Vom, respectively, change. In this example, the current through RL may be calculated as IL=(Vop−Vom)/RL. The change in load current IL is accommodated by a change in the current through the transistors Q1 and Q2 because the DC bias circuit 102 maintains a fixed Ip2 and Im2. As shown in FIG. 2, when Vop−Vom is at a maximum, Ip1 is at a maximum and Im1 is at a minimum. Similarly, when Vop−Vom is at a minimum, Im1 is at a maximum and Ip1 is at a minimum. The current swing of Ip1 and Im1 is offset by Ip2 and Im2.\nThe modulation index for this swing in transistor current may be defined as IQ(max, min)/IQ(DC), where IQ (max, min) is the maximum or minimum total current through the transistor (i.e., bias current +/−load current) and IQ(DC) is the bias current through the transistor, which is substantially constant. Even a small change in load current IL may lead to distortion because the change in load current is reflected by a proportional change in transistor current for the output transistors Q1, Q2. The larger the load current swing, the larger the modulation index, and the larger the distortion. This is because the transconductance of Q1 and Q2 varies as the current through these transistors varies. A variation in transconductance results in a variable gain because the base-emitter voltage drop in Q1 and Q2 varies with changes in transistor current.\nOne way to reduce the modulation index in the driver of FIG. 1 is to increase the DC bias current. By increasing the DC bias, the significance of variations in the load current can be reduced. Although the modulation index can be reduced in this manner, it merely masks the continuing problem of the distortion in the driver 10. Also, reducing distortion by increasing DC bias involves greater power consumption and may be unacceptable in low power applications."}
{"text": "Many working environments present explosion hazards or present a risk of electrical shock. One hazardous environment is a pumping station or a tank that contains hazardous and/or flammable or volatile gases or liquids and chemicals to be pumped. Float switches are commonly used in applications of this sort to detect level for activation of a pump (see FIG. 2), but float switches generally have electrical current that passes through wires and a switch housed within the float, such as a mercury switch located in the float. Wires from the float switch run to a control panel (or other device) located external to the tank or pit and are located outside the hazardous area. Some specialty control panels are explosion proof and can be located in the hazardous area. As liquids rise in the tank, the float tilts and a ball or conductive liquid, such as mercury, moves and makes contact with an electrical switch or contacts of some sort causing the switch to activate. Electrical current then passes from the control panel through the wires, to the switch, completing the circuit. These all present a spark hazard, and if a breakdown in insulation occurs along the electrical path, an explosion can result.\nSome tanks containing flammable liquids or gases use ultrasonic level detection which sends a sonic burst to the surface of the liquid and then back. The transit time of the beam is used to determine the liquid level (some alternative devices use radar or microwave radiation as an energy packet instead of a sound wave, and other sensing technologies are used in level detection, e.g. magnetostrictive, submersible pressure transducers, bubblers, capacitance, etc.). No electrical current is used within the tank or pit and the transmitter and receiver are located external to the hazardous atmosphere, usually mounted on the external tank surface, with the sensing device positioned in the tank. This technique is not generally used to remotely signal a device, such as a signal to engage/disengage a pump (such as a dosing pump) at discrete levels, as these level detectors will detect all fluid levels requiring additional logic circuits to select a predetermined height or level for operation of a pump, thereby raising the complexity and expense of such a system.\nIn some hazardous environments, explosion proof containers are used to contain equipment or devices that may present a possible sparking hazard, such as controls, pumps, motors, etc. While pumps or other devices located in a hazardous environment may be contained in an explosion proof housing, these devices must be activated or deactivated by electrical signals (e.g. providing power to the device). Activation is done remotely from the hazardous environment to reduce the possibility of explosion. Hence, when an operator is onsite, the operator cannot manually activate/deactivate the device within the hazardous area unless the activation device is in an explosion proof housing. It would be desirable to have a switch that could be located within the hazardous environment that could be used to manually activate/deactivate the powered device, and have the switch not present an arcing hazard, and would not have to be located in an explosion proof housing."}
{"text": "The present invention relates to quadrature frequency converters, and more particularly to an automatic quadrature network with phase and amplitude detection to eliminate both phase and amplitude errors in a frequency converter output.\nQuadrature (I/Q) frequency converters rely on having a pair of mixers which are driven by a quadrature local oscillator signal, i.e., two signals of the same frequency, but ninety degrees out of phase. One of the simpler methods of achieving the desired result is to input a single frequency to a quadrature network having two phase shifting paths, one path with a series resistor followed by a shunt capacitance (RC), and the other path with a series capacitance followed by a shunt resistance (CR). If the resistances are equal and the capacitances are equal, then there is a particular frequency for the single frequency input where the outputs of the two paths are of equal amplitude and in quadrature phase. Outside of the particular frequency, while the phase relationship is held, the amplitudes may vary.\nOne method of generating quadrature signals from a local oscillator (LO) uses amplitude detectors and a feedback control loop which adjusts either resistance or capacitance values in the RC and CR paths until the quadrature signals are matched in amplitude. This method is illustrated by U.S. Pat. No. 5,644,260 (DaSilva) and G.B. Patent No. 1,345,274 (Ratzel). However, this method cannot generate perfect phase quadrature and equal amplitudes when the non-varied elements of the network don't match, when the variable elements don't track exactly, or when parasitic elements exist. For example, because of some non-ideal components—an R is too big or too small compared to the other or a C has the same problem, etc.—the frequency of the RC circuit in one quadrature signal path is too low and the frequency of the CR circuit in the other quadrature signal path is too high. FIG. 1 is a graphic representation of the DaSilva implementation, which shows that at the specified frequency—100 MHz—the quadrature signals are equal in amplitude, but the phase differential between the two quadrature signals varies across the frequency spectrum—in this example at the LO frequency of 100 MHz the phase difference is 100° instead of 90°. In order to achieve the proper phase relationship, a static Vcal signal is used in conjunction with the detected amplitudes, the Vcal signal being selected from a table for the frequency of the particular local oscillator signal input. However, the Vcal signal cannot account for dynamic changes to the network due to time and temperature, so the network needs either frequent calibration or always has some phase error present.\nAn alternative method, if amplitude matching is not important, uses a phase detector and feedback loop to adjust one quadrature (RC) path or the other quadrature (CR) path until perfect phase quadrature is achieved, as illustrated by U.S. Pat. No. 4,908,532 (Chadwick). Ignoring amplitude matching, however, may cause amplitude errors, even though the signal is passed through limiting stages. Because the phase difference may be 90° over a wide range of amplitude differences, one side or the other might be so starved of signal that the limiting amplifiers cannot limit. This causes poor match between the two mixers, since the local oscillator level is not the same at the inputs to the mixers.\nWhat is desired is an automatic quadrature network for a quadrature frequency converter that provides equal amplitude and ideal quadrature phase for quadrature signals derived from a local oscillator."}
{"text": "1. Field of the Invention\nThe present invention relates to a device for determining the histogram of a particle-size distribution for products or substances composed of particles, grains or discrete cells and, more especially, to such a device wherein a histogram is provided by analyzing a diffraction spectrum of the substance in a coherent optical system.\nThe device of the invention permits, in particular, a haematological analysis of blood smears or blood particles in liquid flow. It may also be employed in the granulometric study of industrial powders.\n2. Description of the Prior Art\nThe method of analysing powders by diffraction in coherent optical systems was introduced by N. L. Anderson and R. E. Beissner in an article entitled \"Counting and Classifying Small Objects by Far-Field Light Scattering\", Applied Physics, July, 1971, Volume 10, No. 7, pages 1503-1508. This method provides an overall analysis of the sample, giving quick results with sound statistical validity. The principle thereof is as follows: the sample (powder, blood smear, cells in a liquid medium, etc) is illuminated by a laser. By taking measurements from the diffraction spectrum, the characteristics of the population under examination may be deduced.\nMarcel Bessis and Narla Mohandas, Transactions of the Academy of Sciences of Paris, June 17, 1974, Volume 278, pages 3263-3265, and April 26, 1976, Volume 282, pages 1567-1570, were the first to employ analysis by diffraction of biological samples, with a rather particular application in mind: namely, the study of cell deformability.\nTwo pieces of equipment for granulometric analysis by diffraction are on the market: Compagnie Industrielle des Lasers, a French Company of 91460 Marcoussis, France, uses a laser granulometer for the study of powders, by responding to 8 measurement points on the histogram, and Leeds and Northrup Company, North Wales, PA 19454, makes a device referred to as Microtrac which responds to 13 measurement points on the histogram.\nThe foregoing systems and devices process a predetermined number of discrete measurements at a limited number of measurement points. Although this is suitable for industrial powders, the use of a limited number of measurement points is not usually suitable for biomedical applications where a greater resolution in the histogram is called for.\nThe objective of the invention is to provide a continuous histogram of the distribution of particle sizes or, at least a very high number, a thousand or more, of points of this histogram."}
{"text": "The invention relates to a method for the regeneration of acidic, organic lubricant stripped during the cold forming and processed, as well as reaction lubricant residues containing multivalent metal ions and phosphate ions.\nAcidic reaction lubricants have an oily character and are distinguished in that, in addition to an organic lubricant, they contain constituents which are capable of forming a conversion covering on the surface of the workpiece to be transformed U.S. Pat. No. 3,525,651, U.S. Pat. No. 2,739,915, DE-B-21 02 295, EP-A-24 062, EP-A-25 236). The reaction lubricants, which are in practice the most important by far, contain as component which is capable of forming a conversion covering phosphoric acid polyvalent metal ions phosphate. In the treatment of the workpieces to be deformed iron phosphate layers are essentially produced the presence of phosphoric acid, of polyvalent metal phosphates being produced with the use of reaction lubricants containing polyvalent metal phosphates. The reaction lubricants usually applied when dipping generate on the surface of the workpiece a firmly intergrown conversion covering which guarantees that during the deformation a high measure of separation of workpiece and deforming tool takes place, and that on the other hand, a secure binding of the organic lubricant to the workpiece surface is given.\nCold forming happens most frequently with the use of reaction lubricants takes place by means of slide-drawing. In this respect, the excess reaction lubricant is stripped off at the outer drawing matrix.\nIf the slide-drawing takes place with the use of a drawing mandral, excess reaction lubricant is discharged, in addition, from the pipe interior to the outside. The reaction lubricant residues stripped in this way are collected and usually disposed of as waste. A reuse of the stripped reaction lubricant residues is not possible because a change has taken place, as a result of temperature strain and pressure load, such that its original property of layer formation has been completely lost. In addition, rubbed-off metal parts are produced as a result of the friction forces prevailing during the deforming procedure. These are present in undissolved form, but are also present in dissolved form as organic and is inorganic iron compounds, as a result of reaction with further components of the reaction lubricant. The metallic iron content in the reaction lubricant residue prevents its reuse as a result of scratch formation on the pipes to be drawn. In addition, organic iron compounds thicken the reaction lubricant in such a way that a dipping or circulating process can no longer be carried out economically as a result of high lubricant consumption.\nThe aforementioned disposal of the residues of reaction lubricants usually takes place by burning. In this respect, however, not only are combustion products produced, which can represent a considerable strain on the environment, but components of the reaction lubricant which are valuable in themselves are destroyed.\nThe object of the present invention is to provide a method for the regeneration of acid reaction lubricant residues, stripped during the cold forming and processed, which allows the restoration of a functional reaction lubricant with the simultaneous use of a substantial part of the components present in the reaction lubricant residues.\nThe object has been achieved by developing the method of the above-indicated type in accordance with the invention in such a way that, there is added to the reaction lubricant residues, which have passed through the procedural steps of:\na) precipitating metal ions from organometallic compounds by addition of phosphoric acid and/or acid phosphoric acid esters,\nb) dissolving water-soluble impurities by the addition of water, and\nc) separating the lubricant component from solids and from aqueous phases which are incompatible with the oily lubricant component,\nd) acid phosphoric acid ester to adjust the free acid to a value of at least 10 points and to adjust to a total acidity point number of at least 70, with the adjustment of the ratio of free acid: total acidity being to a value in the range of 0.1 to 0.5.\nAccording to the nature of the stripped reaction lubricant residues, it is advisable to add a filtration in order to remove coarse impurities. For this purpose, heating advantageously takes place before the filtration, in order to reduce the viscosity of the reaction lubricant residues.\nAfter the treatment by the addition of phosphoric acid and/or acid phosphoric acid esters corresponding to procedural step a) and the addition of water corresponding to procedural step b), a multi-phase oily product is obtained, the individual phases of which contain the different impurities partially together with the valuable substance to be recovered.\nIn the subsequent procedural step c) the separation of the lubricant component from solids and from aqueous phases incompatible with the oily lubricant component takes place with the aid of filters, decanters, separators, but preferably with the aid of centrifuges. In this respect, all solids and all phases incompatible with the oily phase, for example the aqueous or is aqueous-acid phase, and all phases, the density of which lies above that of the organic lubricant originally used, are separated.\nAccording to a preferred development of the invention the separation corresponding to step c) takes place by means of a centrifuge at 2000 to 30000 g. The use of a centrifuge at 8000 to 15000 g is particularly advantageous.\nAs a result of the separating process, an oily substance is obtained, which contains substantially all organic components of the original reaction lubricant. However, there is a deficiency of the components responsible for the formation of the conversion covering because, as a result of reaction during the drawing procedure, a considerable part of the reactive component of the reaction lubricant has been broken down. To restore the functional reaction lubricant acid phosphoric acid esters are added to the oily lubricant component, obtained after passing through step c), to adjust the free acid to a value of at least 10 points and to adjust the total acidity point number to at least 70. In this respect, it is to be observed that the ratio of free acid: total acidity is adjusted to a value in the range of 0.1 to 0.5. If necessary, it can be advantageous additionally to add phosphoric acid in order to adjust the free acid.\nAcid esters of saturated or unsaturated fatty alcohols, preferably with a chain length of 16 to 22 C atoms, are particularly suitable as phosphoric acid esters.\nThe adjustment of free acid and total acidity takes place according to the following method of determination.\nFor the titration of the free acid about 1 g of the reaction lubricant is weighed exactly into a 300 ml Erlenmeyer flask. 60 ml xylene is added to dissolve the lubricant. This solution is shaken in a 250 ml shaking funnel a total of four times with a mixture of 10 ml demineralized water and 20 ml ethanol. The separated ethanol/water phases are put into a 300 ml Erlenmeyer flask and about 2 ml of an indicator solution, consisting of 0.1% by weight bromocresol green and 99.9% by weight ethanol, are added.\nThereafter, the titration takes place with 0.1 m alcoholic potassium hydroxide solution until there is change in colour from green to blue. The free acid content is given as\nconsumption (ml) potassium hydroxide solutionxc3x975.6 initial weight (g).\nFor the titration of the total acid content, about 1 g of the reaction lubricant is likewise weighed exactly into a 300 ml Erlenmeyer flask and 100 ml of a solvent mixture, consisting of 0.1% by weight alkali blue (indicator), 42.9% by weight xylene, 28.5% by weight ethanol and 28.5% by weight isopropanol, are added. After complete dissolution of the lubricant the titration takes place with 0.1 m alcoholic potassium hydroxide solution until there is change in colour from blue to red-purple. The calculation of the total acid is given by:\nconsumption (ml) potassium hydroxide solutionxc3x975.6 initial weight (g).\nBy suitable selection of the acid phosphoric acid ester and/or the fatty acid, the viscosity of the recovered reaction lubricant should lie between 8 and 1000 mm2/sec at 60xc2x0 C. The preferred range for reaction lubricants to be used in the dipping method lies between 8 and 100 mm2/sec, for circulating products, the preferred range is 50 to 500 mm2/sec."}
{"text": "Semiconductor wafers comprising micro-mechanical components, for example, micro-mirrors, are known. The micro-mechanical components are formed in a membrane layer, and where the micro-mechanical components are micro-mirrors, the micro-mirrors are of thickness similar to the depth of the membrane layer. In order to minimise the mass of the micro-mirrors, and in turn, the drive voltage required for tilting the micro-mirrors, the micro-mirrors should be relatively thin, and typically, of thickness in the range of 3 μm to 10 μm, and preferably, closer to 3 μm for minimising the drive voltage required to tilt the micro-mirrors. Thus, the membrane layer in which the micro-mirrors are to be formed should be of depth similar to the desired thickness of the micro-mirrors, and thus, in the range of 3 μm to 10 μm, and preferably, in the order of 3 μm.\nThe membrane layer is supported on a handle layer, and in general, a buried insulating layer is disposed between the handle layer and the membrane layer. Generally, both the handle layer and the membrane layer are of silicon, and the buried insulating layer is an oxide layer. The oxide layer may be formed on either the handle layer or the membrane layer and laminated to the other of the handle layer and the membrane layer by a suitable bonding process. The buried oxide layer may be either grown, or may be deposited by, for example, chemical vapour deposition. After the membrane layer and the handle layer have been laminated together with the buried oxide layer disposed between the respective layers, the membrane layer is machined to the desired depth, which as mentioned above is typically in the range of 3 μm to 10 μm. The handle layer typically is of depth in the range of 300 μm to 400 μm while the buried oxide layer is of depth typically less than 500 nm.\nAfter lamination and machining of the membrane layer and the handle layer, a photo-resist layer is deposited on the exposed surface of the membrane layer, which is patterned to define the micro-mechanical components. The membrane layer is etched to form a plurality of trenches which extend through the membrane layer to the buried oxide layer for defining the micro-mechanical components. Where the micro-mechanical components are micro-mirrors, two arcuate trenches are typically formed to define each micro-mirror. The two trenches together define the micro-mirror and a pair of tethers. The tethers extend between the membrane layer and the micro-mirror on respective opposite sides of the micro-mirror for supporting the micro-mirror in the membrane layer. The tethers of each micro-mirror define a pivot axis about which the micro-mirror is tiltable. After the membrane layer has been etched to form the micro-mirrors, the handle layer is then etched to form communicating openings to the respective micro-mirrors. The communicating openings are typically formed by through bores which extend through the handle layer and the buried oxide layer. Typically, the through bores are of cross-sectional area greater than the area of the micro-mirrors, and typically, each through bore is of cross-sectional area similar to the area of the micro-mirror plus the width of the trenches on opposite sides of the micro-mirror. At an appropriate time one or both faces of the micro-mirrors are metallised for providing a reflective surface or surfaces on each micro-mirror.\nThe actual etching of the trenches in the membrane layer to the buried oxide layer for defining the micro-mirrors and their tethers is a relatively straightforward procedure. However, once the trenches have been formed in the membrane layer, the portions of the buried oxide layer bridging the trenches are unsupported on the membrane layer side of the buried oxide layer. Additionally, during etching of the through bores through the handle layer, the buried-oxide layer acts as an etch stop layer. Some of the through bores etch faster than others, due to variation of the etch rate across the semiconductor wafer, and thus, the portions of the buried oxide layer at the end of some of the through bores will be exposed to the through bore etch longer than others. The through bore etch has an etching effect on the portions of the buried oxide layer exposed by the through bores, and those portions of the buried oxide layer which are subjected to the through bore etch for a relatively lengthy duration are etched to a relatively thin depth. Additionally, the through bore etch induces stresses in the portions of the buried oxide layer which are exposed by the through bores, and which have been subjected to the through bore etch for relatively lengthy durations.\nIt has been found that the magnitude of the stresses increases as a function of the duration for which the portions of the buried oxide layer are subjected to the through bore etch and the amount by which the buried oxide layer has been thinned. It has been found that when the buried oxide layer has been thinned to a thickness of less than 200 nm, it is no longer able to withstand the induced stresses, and the unsupported portions of the buried oxide layer commence to curl up at the periphery of the through bores, and thus rupture from the buried oxide layer which is buried between the membrane and handle layers. In other words, the portions of the buried oxide layer adjacent the peripheral edge of some of the through bores which are unsupported by virtue of the arcuate trenches having been formed in the membrane layer curl up into the through bore, thus rupturing from the remainder of the buried oxide layer. Since the tethers which are supporting the micro-mirrors are of relatively small transverse cross-section and are attached to the buried oxide layer, the curling up of the peripheral edges of the portions of the buried oxide layer can cause one or both of the tethers to curl up with the buried oxide layer. This, thus, causes the tether or tethers to rupture from the membrane layer, thus leaving the micro-mirror partly or wholly unsupported in the membrane layer.\nEven if the buried oxide layer does not curl up adjacent the periphery of the through bore, thus leaving the tethers intact, it has been found that the stresses induced in the portions of the buried oxide layer at the end of the through bores, which have been exposed to the through bore etch for a relatively lengthy duration, cause the oxide layer to bow. Since the micro-mirrors are relatively thin, in the range of 3 μm to 10 μm and are attached to the buried oxide layer, the bowing of the oxide layer causes similar bowing in the micro-mirrors. Such bowing leaves the micro-mirrors unsuitable for use.\nThe problem of rupturing of the tethers is particularly common in semiconductor wafers where the depth of the membrane layer is at the lower end of the range of 3 μm to 10 μm, and in turn the micro-mirrors or other micro-mechanical components which are of similarly small thickness. Increasing the thickness of the buried oxide layer, while it may overcome the problem of rupturing of the unsupported bridging portions of the oxide layer, results in another problem. It has been found that the stresses induced in portions of thicker buried oxide layers at the end of the through bores which have been exposed to the through bore etch for relatively lengthy durations are such as to cause bowing of the buried oxide layer at the end of the through bore. Since the micro-mirrors or other micro-mechanical components are attached to the buried oxide layer, bowing of the buried oxide layer causes bowing of the micro-mirrors or other micro-mechanical components, which in many cases leads to rupturing of the component, and where the component is a micro-mirror, the micro-mirror may rupture or one or both of its tethers may rupture.\nThe problem of bowing of micro-mirrors resulting from bowing of the buried oxide layer also increases as the depth of the membrane layer, and thus the thickness of the micro-mirror or micro-mechanical component is reduced, and the thickness of the buried oxide layer is increased, since the thinner the micro-mirror or micro-mechanical component is, the more susceptible it is to bowing, resulting from bowing of the buried oxide layer.\nThere is therefore a, need for a method for forming a micro-mirror and for forming other micro-mechanical components in a semiconductor wafer wherein the risk of rupturing or deforming of the micro-mechanical component or a part thereof is minimised.\nThe present invention is directed towards providing such a method, and the invention is also directed towards providing a semiconductor wafer comprising a micro-mechanical component formed therein."}
{"text": "Carbonic acid diester such as dimethyl carbonate are attracting attention as a raw material substituted for phosgene or dimethyl sulfate in the production of isocyanates, polycarbonates and various drugs and agricultural chemicals. Further, the use thereof as an additive to automobile fuel is being considered.\nCarbonic acid diester has long been produced by reacting phosgene with an alcohol. However, the use of phosgene has drawbacks in that phosgene is highly toxic and that the reaction forms hydrochloric acid as a by-product, which causes corrosion of the apparatus. Therefore, processes for producing the carbonic acid diester without the use of phosgene have been developed and are now being executed on industrial scales.\nFor example, Japanese Patent Laid-open Publication No. 54(1979)-24827 discloses a process for producing a carbonic acid diester, in which an alcohol, carbon monoxide and oxygen react in the liquid phase in the presence of a copper halide catalyst. Furthermore, Japanese Patent Laid-open Publication No. 50(1975)-40528, Japanese Patent Publication No. 61(1986)-8816 and Japanese Patent Laid-open Publication Nos. 62(1987)-81356 and 1(1989)-287062 disclose processes for producing the carbonic acid diester in liquid phase, in which copper and palladium halides are used as catalytic components.\nAlthough having the advantage of using no phosgene, the above processes for producing a carbonic acid diester in the liquid phase have drawbacks in that (1) the reaction apparatus is likely to corrode by the action of a halide catalyst which is used as a solution (2), the catalyst activity is deteriorated rapidly by the water formed during the reaction and (3) it is difficult to separate the catalyst dissolved in the reaction product. In particular, it is requisite that the reaction apparatus is constructed of a corrosion resistant high-quality material to thereby prevent the corrosion of the apparatus and any disaster attributed to the corrosion. This inevitably causes the construction cost of the apparatus to be extremely high.\nProcesses for producing carbonic acid diester by vapor phase reaction have been proposed as substitutes for the liquid phase processes (for example, Published Japanese Translation of PCT Patent Applications from Other States, No. 63(1988)-503460, Japanese Patent Laid-open Publication No. 4(1992)-89458 and International Application Publication WO 90/15791).\nIn the vapor phase processes disclosed in these publications, vaporized alcohol, carbon monoxide and oxygen react in the vapor phase in the presence of a catalyst, and the resultant gaseous reaction product is withdrawn from the reactor. The withdrawn gas is cooled to thereby separate it into a condensed liquid and a noncondensable gas and the carbonic acid diester is separated from the liquid. These vapor phase processes are substantially free from the above drawbacks of the liquid phase processes.\nApart from the above, the oxidative carbonylation of an alcohol is an exothermic reaction. The heat of reaction must be removed for maintaining an appropriate reaction temperature. For example, the calorific values are about 71 kcal/mol in both of the following reactions in which carbonic acid diesters are prepared from methanol and ethanol: EQU 2CH.sub.3 OH+CO+1/2O.sub.2 .fwdarw.(CH.sub.3 O).sub.2 CO+H.sub.2 O, and EQU 2C.sub.2 H.sub.5 OH+CO+1/2O.sub.2 .fwdarw.(C.sub.2 H.sub.5 O).sub.2 CO+H.sub.2 O.\nWhen a catalyst consists of, for example, a copper oxychloride supported on an active carbon, the above oxidative carbonylation of an alcohol in the vapor phase is conducted at a temperature as relatively low as about 130 to 170.degree. C.\nWhen it is attempted to synthesize a carbonic acid diester by the use of a fixed-bed reactor, hot spots are likely to occur because the reaction is highly exothermic, so that the dangers of decrease of reaction selectivity, runaway of reaction and catalyst deactivation are increased.\nBy contrast, the heat of reaction can be removed much more easily in a fluidized-bed reactor than in the fixed-bed reactor, so that the temperature can be controlled without the occurrence of hot spots.\nIn the above vapor phase reaction in the fluidized-bed reactor, generally, cooling pipes are inserted in the fluidized bed so as to effect heat removal and water is generally used as the heat transfer medium. For example, International Application Publication WO 95/21692 discloses a fluidized-bed reactor for vapor-phase exothermic reaction which is provided with cooling pipes capable of feeding a cooling medium at a steady rate and cooling pipes capable of feeding a cooling medium at a variable rate. In the publication, water is mentioned as the cooling medium fed at a steady rate and steam is mentioned as the cooling medium fed at a variable rate.\nIf steam obtained by the vaporization resulting from the removal of the heat of reaction generated in the fluidized bed can be used as, for example, a heat source for a reboiler of a distillation column, effective utilization of the heat of reaction would be attained. When the difference between the reaction temperature and the temperature of the heat transfer medium flowing through the cooling pipes is large, the removal of the heat of reaction would be easy.\nIn the oxidative carbonylation conducted at relatively low temperatures, the difference between the reaction temperature and the temperature of steam available as a heat source (for example, about 125.degree. C. with respect to reboiler heat source) is so small that the number of cooling pipes must be increased to thereby enlarge the heat-transfer area in order to obtain the steam of the above temperature.\nHowever, the number of cooling pipes which can be inserted in a fluidized bed of a given volume is so limited that it is likely that the heat-transfer area required for the removal of heat cannot be satisfactorily secured. Further, when the removal of heat is attempted by increasing the difference between the reaction temperature and the temperature of heat transfer medium under the condition of the limited number of cooling pipes, only low-temperature steam or water can be obtained which has little value.\nOn the other hand, the attempt to increase the heat-transfer area leads to an unnecessary expansion of the volume of the fluidized bed.\nTherefore, there has been a demand for the development of a process capable of producing a carbonic acid diester in high energy efficiency with the effective utilization of the heat of reaction in the vapor-phase oxidative carbonylation of an alcohol which is conducted at relatively low temperatures."}
{"text": "Typically, most of the power interruptions to which lamps are subjected are of very brief duration, e.g., under one second, but even an interruption of this brief duration will cause a typical HID lamp to be extinguished. If the usual HID lamp is subjected to such a brief power interruption, it does not immediately reignite upon restoration of normal voltage thereacross. Typically, some cooling of the gases within the HID lamp will be required before the lamp will reignite in response to voltage applied thereacross. When the HID lamp does reignite in response to such voltage, its lumen output at the time of reignition is only a small fraction of normal. Gradually, this lumen output increases until the lamp regains its full brightness."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for optical scanning, and more particularly to a method and apparatus for optical scanning capable of reducing a dot position displacement and a dot diameter variation based on a variation in wavelength of a laser beam.\n2. Discussion of the Background\nBackground optical scanning apparatuses for use in various image forming apparatuses including laserprinters, facsimile machines, and digital copiers have a drawback in which a focal image position on an imaging surface is varied and consequently an image quality on the imaging surface or on an output recording sheet is degraded. This variation in focal image position on the imaging surface is caused when a wavelength of a laser beam generated by a laser light source is varied and a refractive index of optical devices are accordingly changed.\nA technique described in Japanese Laid-Open patent publication, No. 09-021944 attempts to avoid the effects from the variations in wavelength using a dispersion nature of a glass lens to improve an accuracy of output images.\nOn the other hand, an optical scanning apparatus capable of simultaneous multi-beam scanning using a plurality of laser light sources has been developed in response to an increasing trend for a high speed image reproduction. In the multi-beam method, an image quality may be degraded when the multiple beams are not precisely synchronized which is caused in many cases by variations in wavelength of the laser beams.\nA technique described in Japanese Lain-Open patent publication, No. 09-076562, relates to an optical scanning lens to attempt to avoid the above-mentioned problem by reducing a difference in wavelength between multiple beams below 12 nm.\nAlso, in recent years, a demand for a high image density in the image reproduction using an optical scanning apparatus is rapidly increased and, to implement this demand, requirements for reduction of a dot size and a dot position displacement have grown.\nThere is a further problem in that an image quality is degraded due to variations in wavelength of a laser beam caused by changed in temperature. That is, when the wavelength of a laser beam is changed, an optical scan-imaging device such as fθ (ef-theta) lens changes its refractive index and, as a consequence, a start position in the scanning line on an imaging surface is changed which causes a displacement of dot position.\nThe above-mentioned Japanese Lain-Open patent publication, No. 09-021944 proposes a glass-made scanning lens which prevents a deterioration of an image forming accuracy by a reduction of chromatic aberration of the scan-imaging device. However, it is difficult to achieve a similar effect using a plastic lens, which falls into a complex lens structure."}
{"text": "In the course of integrated circuit (IC) evolution, functional density (i.e., the number of interconnected devices per chip area) has generally increased while geometry size (i.e., the smallest component (or line) that can be created using a fabrication process) has decreased. This scaling down process generally provides benefits by increasing production efficiency and lowering associated costs. Such scaling-down also produces a relatively high power dissipation value, which may be addressed by using low power dissipation devices such as complementary metal-oxide-semiconductor (CMOS) devices. During the scaling trend, various materials have been implemented for the gate electrode and gate dielectric for CMOS devices.\nFor example, high-dielectric constant gate dielectric materials and metal gates (HK/MG) have been adopted to enhance electrical performance of CMOS devices. During a gate-last HK/MG process, polysilicon gate electrodes are formed in conjunction with other process for forming the CMOS devices. In a later stage, the polysilicon gate electrodes are removed and replaced with one or more metallic materials."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of data communications for personal computers, and in particular to a system for dynamically collecting and displaying performance improvement information.\n2. Related Art\nThe Internet is quickly becoming the preferred data communications medium for a broad class of computer users ranging from private individuals to large multinational corporations. Such users now routinely employ the Internet to access information, distribute information, correspond electronically, and even conduct personal conferencing. An ever-growing number of individuals, organizations and businesses have established a presence on the Internet through \"Web pages\" on the World-Wide Web (WWW).\nAt the same time that the Internet's popularity has been growing, the trend in personal computer (PC) design has been towards increased portability. This trend is best evidenced by the rapid growth in popularity of so-called \"notebook\" or \"laptop\" computers. Unfortunately, while PCs are becoming both more powerful and more portable, advances in communications infrastructure have not necessarily kept pace.\nA primary limitation on the ability of personal computers to exploit the Internet to full benefit is a lack of sufficient communications bandwidth. Similar problems may be caused by such conditions as high network traffic, even for high-speed Internet access lines. Outside of an office environment, for example, PCs typically connect to the Internet via POTS (Plain Old Telephone Service) or wireless modems at speeds ranging from 9.6 to 14.4 kbps. Such speeds are substantially lower than typical office LAN (Local Area Network) bandwidths. This limited data communications capability can translate into long user-visible latencies when interacting with the Internet.\nA number of approaches have been (and are being) developed to address the problem of slow data communications in the context of Internet applications. One such approach involves using one or more intermediate devices, such as a \"network proxy,\" arranged between a client device and content servers on the Internet to enhance client/network data communications. Network proxies may be used, for example, to cache network content so that subsequent client requests may be satisfied without having to re-retrieve a requested data object from a content server. Network proxies may also be used to \"transcode\" data objects into a form that may be downloaded more quickly. Such transcoding might comprise, for example, techniques such as data compression and image quality reduction which tend to reduce the amount of data to be downloaded.\nAs techniques for improving the performance of client/network data communications are being developed, it is desirable to be able to quantify the performance benefits derived from those techniques and to communicate such information to users. There is therefore a need for methods and apparatus to collect and display performance improvement information for a client device, including when performance improvement benefits are provided to a client device in a distributed manner by a plurality of other devices."}
{"text": "The electrical lighting system of a towing vehicle typically includes brake and turn signal lights mounted on the rear of the vehicle for indicating vehicle braking or directional changes. Specifically, illuminated brake lights indicate a slowing or stopping vehicle, while illuminated and flashing turn signal lights indicate a turning vehicle. Typically in late model vehicles, the turn signal lights and brake lights are independently connected to the electrical lighting system of the towing vehicle and have separate filaments. The brake light filament is typically connected directly to the vehicle brake, while the turn signal light filament is typically connected directly to the turn signal indicator.\nWhen the vehicle is used to tow a trailer, the lights on the vehicle are obscured, so it is necessary to include a separate set of brake and turn signal lights on the trailer. Each brake light and turn signal light on the trailer, however, is typically combined into a single light having a common filament. The single filament functions both as a turning indicator as well as a slowing or stopping indicator.\nConsequently, in order for the trailer light to indicate the correct directional signal, it is necessary to adapt the independent turn signal and brake lights of the towing vehicle to the combined turn signal and brake lights of the trailer. Therefore, a trailer light adaptor circuit is typically inserted between the electrical lighting system of the towing vehicle and the electrical lighting system of the trailer to adapt the towing vehicle's lights to the trailer's lights. The adaptor circuit, for example, can be inserted between the trailer tap terminal on the towing vehicle and the wiring harness of the trailer, as shown in U.S. Pat. No. 4,842,524 issued to the assignee of the present invention.\nPrior art devices have attempted to adapt the electrical lighting system of the towing vehicle to the trailer lights. For example, Bryant, U.S. Pat. No. 4,006,453 shows a device for adapting the independent brake and turn signal lights on a towing vehicle to the combined brake and turn signal lights on a trailer using two logic circuits, one circuit for each side. Bryant discloses logic circuits having npn transistors, silicon controlled rectifiers or relays connected between the electrical lighting system of the towing vehicle and the trailer lights.\nAlthough Bryant discloses a circuit for adapting the electrical lighting system of a vehicle to the trailer lights, there remains a demand for new and improved electrical circuits which perform the same general function while providing manufacturing, operational and cost efficiencies. In this regard, late model vehicles typically have on-board computers and sophisticated electronics to control the engine function and the electrical lighting system in the towing vehicle. The prior art electrical circuits may in some instances permit some feedback from the electrical lighting system of the trailer to the electrical lighting system of the towing vehicle."}
{"text": "1. Field of the Invention\nThe present invention relates to outdoor grills. More particularly, the present invention relates to an outdoor grill of the type that employs a propane or butane fueled burner element placed under a flat grilling surface of a table top, wherein a specially configured tabletop and pedestal arrangement enables a cook to stand next to and cook food on the tabletop grilling surface, and wherein a plurality of individuals are able to sit next to the table at positions next to a cool zone that has quickly removable tiles for easy replacement spaced away from the cook.\n2. General Background of the Invention\nTeppanyaki is a type of Japanese cuisine. It can also refer to a type of grill. In Japan, Teppanyaki can refer to a number of different dishes that are cooked using a flat iron or other metal plate that is heated. It is popular to locate the hot plate at the center of a diner's table. In America, the Benihana® Restaurant chain has for years employed a Teppanyaki style grilling table for the cuisine that they served.\nTeppanyaki style or Asian style table grills are commercially available. One example is a Teppanyaki Asian table grill by Silit-Werke Gmbh & CoKG and sold at www.amazon.com.\nA website that offers a Teppanyaki style cooking apparatus is www.cook-n-dine-usa.com. The product offered by Cook 'N Dine.com provides a cooking center which sinks down slightly and thus forms a shallow pit of up to sixteen (16) inches diameter (thermal expansion towards the heat source underneath). A second area is a warming area, said to be a ring of about three (3) inches width surrounding the cooking center. The temperature in that warming area is approximately fifty percent (50%) lower than the set cooking temperature at the cooking center. A third area is a dining area which is the remaining space toward the edge of the table surface which stays completely cold. The Cook 'N Dine.com product provides a foot having stabilizer arms, a pipe or stand, a heating element and a heating element enclosure with a tabletop cooking surface that includes the three aforementioned zones or areas.\nLarge and relatively expensive, Teppanyaki style grills are available from www.Become.com, such as for example the model DI-48-12KW Drop In Teppanyaki Grill which measures 48″H×23″W×1″deep. A larger 60″×23″×1″ model is also available. Similar commercial model Teppanyaki grills are available from www.instawares.com.\nThe following possibly relevant U.S. patents are incorporated herein by reference:\nTABLEPAT. NO.TITLEISSUE DATE5,740,723 Portable Outdoor High TemperatureApr. 21, 1998Griddle4,724,823Radiant Gas Burner AssemblyFeb. 16, 19886,995,340Portable OvenFeb. 07, 2006"}
{"text": "The present disclosure relates in general to healthcare, and in particular to a system and method for providing healthcare-related services, including, but not limited to, health insurance benefit activation services and/or healthcare-related education services."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device and a backlight module thereof, and more particularly, to a backlight module providing a backlight source to a liquid crystal display panel for generating a two-dimensional/three-dimensional switchable display effect and a display device including the backlight module.\n2. Description of the Prior Art\nThe principle of stereoscopic display technology includes respectively delivering different images to a left eye and a right eye of a viewer, giving the viewer a feeling of gradation and depth in the images, and generating the stereoscopic effect in the cerebrum of the viewer by the analyzing and overlapping of images separately received by the left eye and the right eye.\nIn general, the stereoscopic display technologies may be substantially divided into two major types, which are the glasses type and the naked-eye type (autostereoscopic type). Recently, the parallax barrier type stereoscopic display technology and the lenticular lens type stereoscopic display technology are the popular naked-eye type stereoscopic display technologies. In the parallax barrier type and the lenticular lens type stereoscopic display technologies, stripe barrier patterns or lenticular lenses are disposed in front of a general display panel, and one eye of a viewer may receive images different from those received by another eye of the viewer. The different images received respectively by two eyes of the viewer may be integrated in the viewer's brain to generate stereoscopic display effect. In the conventional naked-eye type stereoscopic display technologies, two adjacent pixels are used to respectively provide different images to the left eye and the right eye at the same time resulting in the perceived resolution becoming lower than the actual resolution of the display panel, and the display quality may then be seriously influenced."}
{"text": "The invention relates to a laser amplifier system comprising a solid body having a laser-active medium, a pump radiation source for producing a pump radiation field which passes multiply through the solid body, and a focusing system which produces a plurality of different branches of the pump radiation field that enter the solid body, and in this case converts at least one branch that emerges from the solid body into one of the branches that enter the solid body and are different from the emerging branch.\nSuch a laser amplifier system is known from the prior art, for example from EP 0 632 551.\nIn these laser amplifier systems, the branch that emerges from the solid body is merely deviated and refocused onto the solid body.\nThe disadvantage of such a focusing system is that, with it, the cross section of the pump radiation field to be focused is always increased or significant losses arise because a part of the pump radiation field is not refocused back onto the solid body.\nOverall, achieving the highest possible pump power radiation density in the solid body is problematic, especially under the premise that the greatest possible number of passes of the pump radiation field is striven for, since the absorption by the laser-active medium in the body during a single pass of the pump power is small.\nIt is therefore an object of the invention to improve a laser amplifier system of the generic type in such a way that the highest possible pump power density in the solid body can be achieved with the least possible outlay."}
{"text": "It is common practice to cover a structure with sheet material for purposes of decoration, protection or reinforcement of the structure. A typical example of a decorative application is wallpaper applied to a wall or similar surface, in order to provide an aesthetically pleasing appearance or a particular desired design. Another typical application utilizing sheet material relates to the construction and design of models, such as scale replicas of airplanes, boats, automobiles or other structures. Such models are usually constructed of a series of spaced apart ribs carved from balsa wood sheets that collectively define a contour such as a camber along a wing or the general shape of the body of the model. A sheet material is applied upon the network of ribs thereby forming a generally smooth continuous surface of the model as defined by the underlying balsa wood structure. The sheet material is preferably comprised of a thin film which serves to define a three dimensional shape and protect the balsa wood structure from exhaust or moisture, as well as to reinforce the balsa wood structure.\nIn such applications, it is highly desirable to cut the sheet material during application in a smooth, straight line. It is also usually highly desirable to cut the sheet material during application at a predetermined uniform distance from a border or edge. A uniform cut likely improves adherence of the sheet material to the structure to which it is applied, and is generally a prized aesthetic characteristic."}
{"text": "1. Field of the Invention\nThe invention pertains to a method and apparatus for displaying light representative of a standard color adjacent a sample, and more particularly to a method and apparatus for displaying light reflected from a standard color adjacent a microscopic sample whereby the color of the sample may be directly compared with the color of the standard color in the field of view of a microscope.\n2. Description of the Prior Art\nFor various reasons, it is often necessary to determine the precise color of a particular sample. This is particularly true of samples being subjected to microscopic examination to classify the sample's origin, composition or background. For example, in the field of palynology, the thermal maturation of spores and pollen that have been preserved in rock formations can be determined by classifying their color. Spores and pollen from relatively young formations are immature and appear to be a light greenish yellow. Spores and pollen from relatively old formations are mature and appear black. Spores and pollen from formations of intermediate age and maturities range in color from yellow through orange to brown. The color of the spores and pollen can be correlated with the thermal history of the formation to give valuable information relating to the potential for producing hydrocarbons from the formation. In analyzing such samples it is critical to distinguish between close shades of colors in order to properly classify the spore or pollen samples as to the degree of thermal maturation.\nColor determination of microscopic samples has commonly been accomplished by microscopic examination of the samples and then subjectively comparing the samples with standard colors on a color chart. Since this process has commonly required the operator to examine the sample in the microscope and then take his eyes away to match the color of the sample with a color on a color chart, a high degree of subjectivity in color determination can result and lead to varying classification, even when the classification is done by the same operator under substantially uniform viewing conditions. The same problem is also present, although somewhat less acute, when non-microscopic samples are visually compared with a color standard to determine the standard color of the sample. Color determination can also be made by spectral examination of light reflected from a sample. Such processes, however, commonly require a relatively large investment in equipment and may be cumbersome in operation."}
{"text": "a) Field of the Invention\nThe present invention is concerned with an improvement to the kind of merchandise display racks known in the trade as \"gondolas\", which improvement essentially lies in the structure and manufacture of the posts and tie-rods acting as structural members in such racks.\nb) Brief Description of the Prior Art\nGondolas are merchandise display racks that are widely used in retail stores especially food stores, to store and display the merchandises offered for sale. Such racks that are most of time disposed in island formation in the stores, basically comprise two or more vertical posts made of rigid metal, such as steel. Each post has a bottom end and a pair of front and rear surfaces formed with a plurality of vertical aligned apertures for use to detachably secure merchandise shelves through angular brackets in an overhanging fashion.\nAt least two and preferably more than two horizontal tie bars also made of metal, are used for rigidly interconnecting each pair of posts adjacent each other in spaced apart relationship.\nThe posts interconnected by the tie-bars are mounted onto the floor by means of transversal footings connected to their bottom ends and sized and positioned to hold the posts vertical even when shelves are secured thereto and loaded with merchandises to be displayed. Each footing is defined by at least one half-base having a vertical inner edge from which hooks projects, and a vertical outer edge. Each half-base is rigidly connectable to the post adjacent the bottom end thereof by insertion of its hooks into the apertures made in the front or rear surfaces of this post. When two half-bases are connected to one post, they horizontally projects away in opposite directions from the front and rear surfaces of the post and act as symmetrical stabilizers for holding this post vertical.\nThis basic structure is quite efficient. However, it calls for the manufacture of a plurality of structural members of different shape, namely the posts and the tie-bars which usually are of different structure depending on when they are intended to be positioned along the posts to interconnect the same. This in turn calls for different manufacturing processes (cold-forming, punching, welding . . .) which are sometimes difficult to carry out in line in a rational, time-and-cost efficient manner."}
{"text": "The IEEE standard, “IEEE Std 1394-1995 Standard For A High Performance Serial Bus,” is an international standard for implementing an inexpensive high-speed serial bus architecture which supports both asynchronous and isochronous format data transfers. Isochronous data transfers are real-time transfers which take place such that the time intervals between significant instances have the same duration at both the transmitting and receiving applications. Each packet of data transferred isochronously is transferred in its own time period. An example of an ideal application for the transfer of data isochronously would be from a video recorder to a television set. The video recorder records images and sounds and saves the data in discrete chunks or packets. The video recorder then transfers each packet, representing the image and sound recorded over a limited time period, during that time period, for display by the television set. The IEEE 1394-1995 serial bus architecture provides multiple channels for isochronous data transfer between applications. A six bit channel number is broadcast with the data to ensure reception by the appropriate application. This allows multiple applications to simultaneously transmit isochronous data across the bus structure. Asynchronous transfers are traditional data transfer operations which take place as soon as possible and transfer an amount of data from a source to a destination.\nThe IEEE 1394-1995 standard provides a high-speed serial bus for interconnecting digital devices thereby providing a universal I/O connection. The IEEE 1394-1995 standard defines a digital interface for the applications thereby eliminating the need for an application to convert digital data to analog data before it is transmitted across the bus. Correspondingly, a receiving application will receive digital data from the bus, not analog data, and will therefore not be required to convert analog data to digital data. The cable required by the IEEE 1394-1995 standard is very thin in size compared to other bulkier cables used to connect such devices. Devices can be added and removed from an IEEE 1394-1995 bus while the bus is active. If a device is so added or removed the bus will then automatically reconfigure itself for transmitting data between the then existing nodes. A node is considered a logical entity with a unique address on the bus structure. Each node provides an identification ROM, a standardized set of control registers and its own address space. Because of these advantages, the IEEE 1394-1995 standard provides for a unique networking structure that is capable of incorporating audio/video devices, media play/record devices and computing/display devices.\nA diverse range of products can be implemented with the ability to connect to an IEEE 1394-1995 serial bus network. These devices can have capabilities and functionality ranging from very simple to very complex. Specifically, a variety of audio/video devices, media play/record devices and computing/display devices are capable of being linked together over an IEEE 1394-1995 serial bus networking structure to support asynchronous and isochronous data transfers between the devices.\nThe IEEE 1394-1995 serial bus allows a collection of devices to work together in a high bandwidth, distributed environment to maximize the overall efficiency and functionality of the network. This allows manufacturers to remove expensive pieces of functionality from one device and locate that functionality in another device on the network, instead of duplicating this functionality in all devices on the network. While some of the devices have limited functionality and are relatively inexpensive, such devices require the support and interaction of other devices in order to bring the full functionality of the devices within the network to the user.\nThe AV/C Digital Interface Command Set is a command set used for data transactions between consumer audio/video equipment over an IEEE 1394-1995 serial bus. Neither the IEEE 1394-1995 serial bus nor the AV/C command set provide a master-slave relationship between the devices coupled within the IEEE 1394-1995 serial bus network. Instead, both the IEEE 1394-1995 serial bus and the AV/C command set operate based on a cooperative peer-to-peer coexistence of devices within the network. Discrete AV/C command and response data packets are transferred between networked devices over an IEEE 1394-1995 serial bus in an asynchronous data stream. The AV/C command and response data packets are typically formatted according to the AV/C protocol outlined in the AV/C Digital Interface Command Set. Transfers of AV/C command and response data packets over the IEEE 1394-1995 serial bus network are supported by an AV/C architecture. The AV/C architecture is comprised of lists and tables that help devices create, process and/or transmit AV/C command and response data packets. The AV/C architecture includes an AV/C bulletin board subunit that is typically dedicated to a resource device subunit such as a tuner, receiver or recoding unit.\nAn AV/C bulletin board subunit is an information structure that is shared between devices networked over an IEEE 1394-1995 serial bus network. A resource schedule bulletin board is also an information structure that supports information shared between coupled devices within a network. The resource schedule bulletin board provides the organizational structure through which shared data is organized and communicated. The resource schedule bulletin board contains sub-boards of lists with entry descriptors that represent encoded data that can be shared between devices within the network via descriptor commands. A dedicated AV/C bulletin board subunit typically supports the information architecture between that device and all compatible posting devices within an IEEE 1394-1995 serial bus network. A posting device writes a request entry to a write enabled list within the resource schedule bulletin board specifying when it will use the resource.\nNeither the IEEE 1394-1995 serial bus nor the AV/C Command Set provide a master-slave relationship between the devices coupled within the IEEE 1394-1995 serial bus network. Instead, both the IEEE 1394-1995 serial bus and the AV/C Command Set operate based on the cooperative peer-to-peer coexistence of devices within the network transmitting data formatted in accordance with the AV/C protocol. The communicating devices exchange command and response data directly with each other without the intervention of a systems resource manager. Each device or device subunit is responsible for managing scheduling affairs and storing resource requests.\nWithin the current AV/C protocol, resource time of a resource device is made available to any requesting control device. Within this protocol a device competing for use of a resource device can acquire the right to use the resource device by reserving the resource device, even though another client device may have previously submitted an entry to the resource schedule board requesting the resource device at the same time. Devices may modify or delete data contained within the schedule entries submitted by other requesting control devices without acknowledging to the other requesting control devices that their respective schedule entries are being or have been modified. Thus, within the current AV/C protocol there is no method to allow control devices within a network to monitor schedule data or to provide a peer-to-peer review of schedule modifications."}
{"text": "Loose-leaf paper binders and covers utilizing binder posts inserted into prepunched paper sheets have been employed for many years. For example, U.S. Pat. Nos. 243,730,560, 3,834,739 and 4,072,326 show plastic strip retainers placed at paper stack edges and integral studs in one retainer strip passed through the punched paper holes. An apparatus is provided to shear and upset the stud heads in a recess of the other retainer strip to hold the stack of punched paper sheets together. Provision is also made to fold a cover around the front and back of the book or report, including a spine portion. U.S. Pat. No. 1,841,989 shows strips of sheet metal which press against a fastener head and surround a ratchet post and wherein collars and wedges are hammered on the posts to compress the sheet metal strips and bind sheets of paper and a cover. Covers which are clamped by a series of fasteners are seen in U.S. Pat. No. 3,667,854. Self adhesive layers with removal strips to expose the adhesive have been utilized in the cover art as shown in U.S. Pat. No. 3,833,244. Wraparound covers in which one edge of the cover overlaps a fastener strip is disclosed in U.S. Pat. No. 4,139,216."}
{"text": "1. Field of the Invention\nThe present invention relates to a rotor employed in a flywheel energy storage device. More specifically, the present invention relates to a rotor wherein the power density and the energy density of the system are optimized while the manufacturing cost is minimized. Methods for manufacturing a rotor according to the present invention are also disclosed.\n2. Description of the Prior Art\nFlywheel systems currently being designed for mobile energy storage are generally intended to replace batteries in electrically powered vehicles. One specific example of such a system was disclosed in commonly assigned, co-pending application Ser. No. 08/148,361, which is incorporated herein by reference for all purposes. In such applications, multiple units are needed to store the required energy, so that each motor-generator need supply only a small portion of the vehicle's power. In systems where all of the surge power must be supplied by a single flywheel, the relatively large size of the single motor-generator makes it difficult to provide the needed energy density without reducing safety factors, e.g., for radial stresses, to unacceptable low levels or raising manufacturing costs to exorbitantly high levels.\nU.S. Pat. No. 3,741,034 discloses rotor designs using high strength-to-weight ratio filament wound composites in relatively thin concentric cylinders, which cylinders are separated by radial springs. While this arrangement limits the radial stresses to tolerable values, it is expensive to manufacture.\nU.S. Pat. No. 3,859,868 discloses techniques for varying the elasticity-density ratio of the rotor elements to minimize radial stresses. On the other hand, U.S. Pat. Nos. 4,341,001 and 4,821,599 describe the use of curved metallic hubs to connect the energy storage elements to the axle. Additionally, U.S. Pat. No. 5,124,605 discloses a flywheel system employing counter-rotating flywheels, each of which includes a hub, a rim and a plurality of tubular assemblies disposed parallel to the hub axis for connecting the hub to the rim while allowing for differential radial expansion between the hub and the rim.\nNone of these references deal with the integration of a large, high power motor-generator into the flywheel energy storage system currently being designed for vehicles."}
{"text": "1. Technical Field of the Invention\nThis invention relates generally to computing systems and more particularly to data storage solutions within such computing systems.\n2. Description of Related Art\nComputers are known to communicate, process, and store data. Such computers range from wireless smart phones to data centers that support millions of web searches, stock trades, or on-line purchases every day. In general, a computing system generates data and/or manipulates data from one form into another. For instance, an image sensor of the computing system generates raw picture data and, using an image compression program (e.g., JPEG, MPEG, etc.), the computing system manipulates the raw picture data into a standardized compressed image.\nWith continued advances in processing speed and communication speed, computers are capable of processing real time multimedia data for applications ranging from simple voice communications to streaming high definition video. As such, general-purpose information appliances are replacing purpose-built communications devices (e.g., a telephone). For example, smart phones can support telephony communications but they are also capable of text messaging and accessing the internet to perform functions including email, web browsing, remote applications access, and media communications (e.g., telephony voice, image transfer, music files, video files, real time video streaming. etc.).\nEach type of computer is constructed and operates in accordance with one or more communication, processing, and storage standards. As a result of standardization and with advances in technology, more and more information content is being converted into digital formats. For example, more digital cameras are now being sold than film cameras, thus producing more digital pictures. As another example, web-based programming is becoming an alternative to over the air television broadcasts and/or cable broadcasts. As further examples, papers, books, video entertainment, home video, etc. are now being stored digitally, which increases the demand on the storage function of computers.\nA typical computer storage system includes one or more memory devices aligned with the needs of the various operational aspects of the computer's processing and communication functions. Generally, the immediacy of access dictates what type of memory device is used. For example, random access memory (RAM) memory can be accessed in any random order with a constant response time, thus it is typically used for cache memory and main memory. By contrast, memory device technologies that require physical movement such as magnetic disks, tapes, and optical discs, have a variable response time as the physical movement can take longer than the data transfer, thus they are typically used for secondary memory (e.g., hard drive, backup memory, etc.).\nA computer's storage system will be compliant with one or more computer storage standards that include, but are not limited to, network file system (NFS), flash file system (FFS), disk file system (DFS), small computer system interface (SCSI), internet small computer system interface (iSCSI), file transfer protocol (FTP), and web-based distributed authoring and versioning (WebDAV). These standards specify the data storage format (e.g., files, data objects, data blocks, directories, etc.) and interfacing between the computer's processing function and its storage system, which is a primary function of the computer's memory controller.\nDespite the standardization of the computer and its storage system, memory devices fail; especially commercial grade memory devices that utilize technologies incorporating physical movement (e.g., a disc drive). For example, it is fairly common for a disc drive to routinely suffer from bit level corruption and to completely fail after three years of use. One solution is to utilize a higher-grade disc drive, which adds significant cost to a computer.\nAnother solution is to utilize multiple levels of redundant disc drives to replicate the data into two or more copies. One such redundant drive approach is called redundant array of independent discs (RAID). In a RAID device, a RAID controller adds parity data to the original data before storing it across the array. The parity data is calculated from the original data such that the failure of a disc will not result in the loss of the original data. For example, RAID 5 uses three discs to protect data from the failure of a single disc. The parity data, and associated redundancy overhead data, reduces the storage capacity of three independent discs by one third (e.g., n−1=capacity). RAID 6 can recover from a loss of two discs and requires a minimum of four discs with a storage capacity of n−2.\nWhile RAID addresses the memory device failure issue, it is not without its own failure issues that affect its effectiveness, efficiency and security. For instance, as more discs are added to the array, the probability of a disc failure increases, which increases the demand for maintenance. For example, when a disc fails, it needs to be manually replaced before another disc fails and the data stored in the RAID device is lost. To reduce the risk of data loss, data on a RAID device is typically copied on to one or more other RAID devices. While this addresses the loss of data issue, it raises a security issue since multiple copies of data are available, which increases the chances of unauthorized access. Further, as the amount of data being stored grows, the overhead of RAID devices becomes a non-trivial efficiency issue."}
{"text": "Stereoscopic, or three-dimensional, imagery is based on the principle of human vision. Two separate detectors detect the same object or objects in a scene from slightly different positions and/or angles and project them onto two planes. The resulting images are transferred to a processor which combines them and gives the perception of the third dimension, i.e. depth, to a scene.\nMany techniques of viewing stereoscopic images have been developed and include the use of colored or polarizing filters to separate the two images, temporal selection by successive transmission of images using a shutter arrangement, or physical separation of the images in the viewer and projecting them separately to each eye. In addition, display devices have been developed recently that are well-suited for displaying stereoscopic images. For example, such display devices include digital still cameras, personal computers, digital picture frames, set-top boxes, high-definition televisions (HDTVs), and the like.\nThe use of digital image capture devices, such as digital still cameras, digital camcorders (or video cameras), and phones with built-in cameras, for use in capturing digital images has become widespread and popular. Because images captured using these devices are in a digital format, the images can be easily distributed and edited. For example, the digital images can be easily distributed over networks, such as the Internet. In addition, the digital images can be edited by use of suitable software on the image capture device or a personal computer.\nDigital images captured using conventional image capture devices are two-dimensional. It is desirable to provide methods and systems for using conventional devices for generating three-dimensional images. In addition, it is desirable to provide systems for providing improved techniques for processing images for generating three-dimensional images or two-dimensional images."}
{"text": "1. Field of the Invention\nThis invention relates generally to a novelty pillow, and more particularly, to a novelty pillow having at least one slit and at least one elongated member, such as a phallic-like member, extending outward from the slit and indicia for conveying a message related to the slit and member.\n2. Description of the Background Art\nVarious types of pillows are known in the art. For instance, there are bed pillows, throw pillows, body pillows, decorative pillows and holiday pillows, to name a few. In fact, numerous patents have issued for various types of pillows. U.S. Pat. Des. 237,569, issued to De Rue, discloses a pillow in the shape of female breasts. U.S. Pat. No. 4,794,657, issued to Avery, discloses a tummy hugger pillow combination comprising trunk, head and leg pillows. U.S. Pat. Des. 391,112, issued to Houston, discloses another pillow design. U.S. Pat. Des. 320,522, issued to Deremiah discloses a combined pillowcase and multiple pillow assembly. U.S. Pat. No. 5,809,594, issued to Isogai, discloses a pillow combination for providing cervical vertebras support. U.S. Pat. Des. 397,901 discloses a pillow displaying a happy face. U.S. Pat. Des. 320,522, issued to Deremiah, discloses a combined pillowcase and multiple pillow assembly. U.S. Pat. No. 5,168,590, issued to O'Sullivan, discloses a therapeutic pillow cover having compartments for receiving hot/cold packs and/or insert supports. U.S. Pat. No. 5,864,904, issued to Rudick, discloses a sectioned bed pillow connected by a fabric tether. U.S. Pat. No. 5,864,903, issued to Newman, discloses a pillow having a storage pocket for storing a condom. Many other patents have issued for various other pillow designs. Despite the number of patents that have issued in this art, none are known which disclose a humorous pillow design having components which resemble a penis.\nA pillow design with features that resemble a penis would be welcomed by the novelty items business and those looking for a new, humorous gift to give at bachelor parties, wedding showers or other events. A pillow with a penis member could be creatively designed with indicia for conveying funny messages. Such a pillow would provide a new and welcomed product to the gag-gift industry. As there are no pillows known having these and similar design features, there exist a need for a penis pillow in the gag-gift industry. The instant invention addresses this need by providing a pillow having at least one elongated member and corresponding indicia."}
{"text": "Though numerous applications utilize coiled heat exchangers, the gas turbine recuperator is among the most demanding. In any application, and especially when used as a gas turbine recuperator, the heat exchanger should be compact, efficient, reliable, and relatively inexpensive to manufacture. By designing the primary heat transfer surface with small hydraulic diameters and counterflow circulation of heat transfer fluids, a relatively compact and efficient heat exchanger can be obtained. Furthermore, providing the heat exchanger with relatively large cross-sectional flow areas reduces load losses. Achievement of large cross-sectional flow areas in coiled heat exchangers requires circulating heat transfer fluids in the axial, as opposed to tangential, direction. Additionally, production costs can be lowered by minimizing the number of elements used to make the heat exchanger and by forming and coiling the heat exchanger in a continuous process. Another design consideration, especially when used as a gas turbine recuperator, includes resistance to thermal shock. Heavy thermal loads often result from the transient operation of turbines. Therefore, to ensure reliable performance and operation, the heat exchanger should have high resistance to thermal shock.\nVarious known heat exchangers are made from coiling a pair of sheets between which heat transfer fluids circulate in a counterflow manner in directions substantially parallel to the longitudinal axis of the coil. For example, U.S. Pat. No. 5,797,449 pertains to an annular heat exchanger formed by a pair of sheets welded together and coiled, with openings cut through the sheets through which heat transfer fluid passes.\nGerman patents DE 1121090 and DE 3234878 describe spiral heat exchangers having axially circulated fluid flows, in which the fluids enter and exit through alternating angular sectors. In DE 1121090, sectors for circulating the heat transfer fluids are formed by cutting evenly-spaced openings in borders that close the edges of a pair of sheets coiled to form the heat exchanger. After the borders are cut, the two sheets are coiled to form the heat exchanger. DE 1121090 additionally discloses the fabrication of the spiral heat exchangers with external headers.\nIn DE 3234878, the sectors are formed by glueing blocking segments on the two faces of the coiled heat exchanger.\nFinally, in French patent document FR-A-2319868, borders are closed by the direct welding of adjacent sheets.\nA particular difficulty in heat exchangers having a coiled configuration includes the distribution of the single incoming flow into the myriad of small heat transfer passages and the collection of the same into a single outgoing flow after the heat transfer has taken place. Preferably, this distributing and collecting should not result in excessive head losses, nor should it cause mechanical stresses due to large thermal gradients. Another difficulty arises from blockages to the heat transfer fluids that exist on the core face as the result of the particular construction used for the heat exchanger. For instance, in one known example, the sheets are constructed and cut such that one sheet has openings only for one fluid and the other has openings only for the other fluid. This leads to a relatively high amount of fluid being blocked at the core faces, thus reducing gas flow passage and overall efficiency of the heat exchanger.\nStacked plate heat exchangers often include openings cut in the plates to distribute and collect the heat transfer fluids. The edges of these openings generally are either brazed or welded together during assembly of the heat exchanger (for example in U.S. Pat. No. 4,073,340) or are fitted with a gasket (for example in Alfa-Laval plate heat exchangers). Other stacked plate heat exchangers do not include such openings (see SAE 851254: \"Development, Fabrication, and Application of a Primary Surface Gas Turbine Recuperator\", E.L. Parsons), but the sides of the plates must be provided with sealing bars."}
{"text": "This disclosure relates to a lubrication system, and more particularly, to a lubrication system for a fan drive gear system in gas turbine engines.\nIn many gas turbine engines, a low pressure spool includes a low pressure turbine connected to and driving a low pressure compressor, and a high pressure spool includes a high pressure turbine connected to and driving a high pressure compressor. A main pump is typically driven by the high pressure spool, connected through gearing, and is used to pump lubricating and cooling fluid to all engine components that require lubrication and cooling.\nThe main pump typically pumps fluid from a passage connected to a main reservoir that holds both liquid and air. During normal operating conditions, the fluid settles at the bottom of the main reservoir and displaces air to the top. However, in a gas turbine engine mounted on an aircraft, the main reservoir may experience “negative gravity” conditions such as the aircraft turning upside down, the aircraft accelerating toward the Earth at a rate equal to or greater than the rate of gravity, or the aircraft decelerating at the end of a vertical ascent. Under negative gravity conditions, the fluid in the main reservoir can rise to the top, which can expose an opening of the passage to air and interrupt the supply of fluid to the main pump and, consequently, interrupt supply to the engine components. Certain engine components, such as gears and bearings, can be damaged by a relatively short period of non-lubricated operation during negative gravity conditions.\nIn some gas turbine engines, a fan at the front of the engine is connected to the low pressure spool through a fan drive gear system. When the high pressure spool stops rotating or rotates at a reduced rpm (revolutions per minute), the fan drive gear system can continue rotating even though the main pump will ordinarily provide little or no fluid during this time. For example, wind may rotate the fan and corresponding gears and bearings while the aircraft is parked on the ground or during an in-flight engine shutdown. Certain gears and bearings can also be damaged by a relatively short period of non-lubricated operation during windmilling as well."}
{"text": "Structured wiring enclosures or “cans” are typically recessed in walls. Known enclosures have limited capability and versatility with regard to mounting electronic devices and/or cable or wire management.\nEnclosure door covers or enclosure doors typically mount to an underlying enclosure. Known covers or doors for structured wiring enclosures do not provide substantial functionality beyond covering and providing access to the enclosure."}
{"text": "The present invention relates to a coding method and a coder effectively using pattern matching, and more particularly to a coding method and a coder using pattern matching capable of coding image data, including character areas and image areas, at a high compression rate.\nA coding method using pattern matching is known as a method for compressing binary data including character areas and image areas. For example, the pattern matching coding method described in ISO/IEC JTC1/SC29/WG1 N205 xe2x80x9cLoseless and Lossy Compression of Text Images by Soft Pattern Matchingxe2x80x9d (issued on Jun. 30, 1995) extracts a small unit (hereinafter called xe2x80x9cmarkxe2x80x9d), for example, one character unit at a time, from the image, and compares it with a plurality of marks stored previously in the library. When the library contains a mark matching the extracted mark, the extracted mark is coded while referencing the mark in the library. The xe2x80x9cmatchxe2x80x9d in this description includes a perfect match as well as a close match (i.e., close to a perfect match).\nFIG. 14 is a block diagram showing the configuration of a coder which executes the coding method described in the above publication. As shown in FIG. 14, a character/image separation unit 200 separates image data received from an input terminal 1 into two types of areas: a character area where characters concentrate and an image area where characters do not concentrate. The image data in the character area is output to a mark extraction unit 201, while the image data in the image area is output to a coding unit 202.\nThe mark extraction unit 200 extracts characters from a character area, one character at a time. The extracted character is treated as a mark. Extracted marks are output to a match checking unit 203 and a storage unit 204. The match checking unit 203 checks if the library contains a mark that matches an extracted mark. If there is a matching mark, the match checking unit 203 outputs the extracted mark and the matching mark found in the library to a coding unit 205. The storage unit 204, which creates the library, adds received marks to the library.\nThe coding unit 205 performs arithmetic coding using the templates shown in FIGS. 15(A) and 15(B). That is, using the reference pixels (indicated by numbers 1-4) in the template, shown in FIG. 15(A), for the extracted mark and the reference pixels (indicated by numbers 5-11) in the template, shown in FIG. 15(B), for the extracted mark , the coding unit 205 estimates the probability at which the pixel level of the processing target pixel is, for example, 1. In general, the coding unit 205 assigns a short code to the probability of 100% or 0%, and a longer code as the probability approaches 50%.\nThe reference pixel position at number 7 in FIG. 15(B) corresponds to the position of P which is the processing target pixel. In the following description, the templates shown in FIGS. 15(A) and 15(B) are called match templates.\nWhen the library does not contain a mark that matches the extracted mark, the coding unit 205 performs arithmetic coding using the template shown in FIG. 16. That is, using the reference pixels (indicated by numbers 1-10) in the template shown in FIG. 16, the coding unit 205 estimates the probability that the pixel level of the processing target pixel P becomes xe2x80x9c1xe2x80x9d and performs arithmetic coding for the estimated probability. In the following description, the template shown in FIG. 16 is called xe2x80x9cunmatched templatexe2x80x9d.\nThe coding unit 202 receives image data of an image area and codes it using a coding method suitable for compressing general binary data, for example, using the JBIG method. Coded data output from the coding units 202 and 205 is output from an output terminal 3.\nAs described above, when the library contains a mark that matches an extracted mark, the match template is used. That is, not only the preceding pixels in the extracted mark (past pixels in the scan direction: pixels 1-4) but also those pixels ahead of (beyond) a position of the processing target pixel (future pixels in the scan direction: pixels 8-11) in the mark having a high correlation with the extracted mark are used to estimate the pixel level of the processing target pixel position P. Therefore, the estimated result is more likely to be 100% or 0% and, as a result, the amount of data is reduced.\nHowever, in the above-described coding method, the coding efficiency of an image area is not high because a standard coding method is used for it. Thus, when the image area is relatively large, the compression rate is not increased. In addition, the need to install two coding units, 202 and 205, increases the size of the coder circuit.\nThe inventor of the present invention has disclosed in Japanese Patent Application No. Hei 10-110835 a coder which checks the dot frequency interval (herein may be termed as simply xe2x80x9cfrequencyxe2x80x9d) in an image area to allow a mark in the library matching an extracted mark to be used even in the image area, that is, a coder which may use a match template. Such a unit, when used, for example, in a newspaper telephotography system having relatively large dot areas, increases coding efficiency and decreases the circuit size.\nHowever, in a system such as a color paper telephotography system, the dot angle of image data in one color differs from that of image data in another color to prevent such occurrence as Moirxc3xa9. Therefore, the dot angle of image data is not always 0 degree or 45 degrees. For example, such a combination is used of a dot angle of 15 degrees in a cyan image, a dot angle of 45 degrees in a magenta image, a dot angle of 0 degree in a yellow image, and a dot angle of 75 degrees in a black image. Note that the dot angle is defined as an angle of a sequence of dots with respect to the sub scan direction of the image.\nIt is easy to find a dot frequency (interval) when the dot angle is 0 degree or 45 degrees. However, it is difficult to find a dot frequency when the dot angle is not 0 degree or 45 degrees. This results in more cases where no mark that matches the extracted mark is found in the library. In this case, the coding efficiency is lower than a case where the library contains the match matching the extracted mark.\nIn view of the foregoing, it is an object of the present invention to provide a coding method and a coder which use pattern matching that ensures high coding efficiency even when the dot angle is not 0 degree or 45 degrees as in a color paper telephotography system.\nAccording to an aspect of the present invention, there is provided a coding method using pattern matching comprises the steps of: separating an image area from the image data; converting a dot angle of the image data in the image area to a specific angle; finding a dot frequency (interval) of the image data whose dot angle has been converted; and dividing the image data, whose dot angle has been converted, by the dot frequency or by an integral multiple of the dot frequency into a plurality of marks. This configuration allows even an image area to be coded through pattern matching and, at the same time, the dot frequency to be found easily, thus enabling marks to be generated effectively using pattern matching.\nThe specific angle for use in dot angle conversion is, for example, 0 degree or 45 degrees so that the dot frequency is easily found.\nWhen the library contains the matching first mark, pixels positioned ahead of (beyond) a position of a processing target pixel may also be included in the reference pixels. The pixels positioned ahead of the position of the processing target pixel are future pixels that are in the scan direction. Therefore, the estimated result is more likely to be 100% or 0% and, as a result, the amount of code data is reduced.\nWhen the library does not contain the matching first mark, the pixels of the mark that are in a predetermined template may be used as reference pixels.\nA newly found mark may be added to the library. This configuration allows the library to grow as coding processing proceeds, ensuring a higher coding efficiency.\nAccording to a second aspect of the present invention, there is provided a coder using pattern matching comprising: a character/image separation unit which extracts an image area from the image data; a dot angle conversion unit which converts a dot angle of the image data in the image area into a specific angle, the image area being extracted by the character/image separation unit; a dot frequency (interval) detecting unit which finds a dot frequency from the image data converted by the dot angle conversion unit; and a mark re-dividing unit which divides the image data, converted by the dot angle conversion unit, by a result found by the dot frequency detecting unit to generate marks for output to the match determination unit. This configuration allows even an image area to be coded through pattern matching and, at the same time, the dot frequency (interval) to be found easily, thus enabling marks to be generated effectively using pattern matching.\nThe mark re-dividing unit, for example, is configured to divide the image data by the dot frequency or by an integral multiple of the dot frequency, the image data being converted by the dot angle conversion unit.\nThe dot frequency detecting unit, for example, is configured to find correlation of the image data converted by the dot angle conversion unit.\nThe dot angle conversion unit may have a plurality of memories containing the image data converted by the dot angle conversion unit and converts the image data by writing the image data into, or reading the image data from, the memories in a direction inclined by a difference angle between the dot angle and the specific angle.\nThe coding unit may comprise a match template processing unit which, if the library contains the matching first mark, uses the pixels of a template of the first mark in the library and pixels of a predetermined template of the second mark in the input image data as reference pixels; an unmatch template processing unit which, if the library does not contain the matching first mark, uses pixels of a predetermined unmatch template of the second mark in the input image data as the reference pixels; and an arithmetic coding unit which estimates a probability of a pixel level of a processing target pixel using reference pixel values from the match template processing unit and the unmatch template processing unit and arithmetically codes the estimation value."}
{"text": "The present invention relates to computing the length of a railway vehicle or a train of such vehicles.\nThe concept of replacing conventional railway signalling equipment by train carried, point controlling, and interlocking modules inter-connected by radio or other transmission medium is well known. There are however certain practical difficulties, missing elements, and unnecessary complications in the realization of known systems. One such is the method of determining the location of the tail (or rear) of a train for the purposes of authorizing a following train to proceed or of allowing a route which has been traversed to be released.\nIn some known systems, train location is determined by an on-board train computer by a combination of reference wayside markers such as transponders, loop ends, loop transpositions or the like together with on-board distance measurement derived from an odometer, tachometer or equivalent device such as a reader of closely spaced track-side marks, plates or loop transpositions or the like. The location thus derived is reported to a wayside control computer which uses the information as the basis for generating movement authorities for transmission to trains with reference also to supervisory controls and fixed interlocking data. For the purpose of route release behind a train and to define the limit of proceed authority for a following train it is necessary for the control computer to deduce the location of the rear of the train. In a known system (British Patent Specification No. 2 189 066)/the control computer does this by reference to the \"train consist\" (i.e. the overall composition of the train) from which train length can be derived. In practice there are difficulties in ensuring that such information is in a sufficiently vital form for use in safety functions. It would be possible to provide equipment at the rear of the train equivalent to that at the front and to report the location of the rear of the train separately from the front, but this would be complicated and expensive and would double the number of mobile identities to be serviced by the communications system. It is also frequently impractical to locate appropriate equipment at the rear of a train."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical displacement meter, an optical displacement measuring method, an optical displacement measurement program, a computer-readable recording medium, and a device that records the program, for irradiating an object to be measured with light, receiving the light from the object to be measured with a light receiving device to obtaining an electric signal according to the receiving light, measuring a distance from a light projector to the object to be measured, a displacement of the object to be measured, and the like corresponding to the electrical signal.\n2. Description of the Related Art\nTo measure the dimension, a movement amount, and the like of an object to be measured (work), an optical displacement meter to which triangulation measuring technique is applied is used. FIG. 93 is a block diagram showing the configuration of main components of a conventional optical displacement meter. In FIG. 93, a drive circuit 101 drives a laser diode (LD) 102 on the basis of a light output control signal Va. A laser beam emitted from the laser diode 102 is passed to a work WK via a projection lens 103. A diffusion reflection component and a specular reflection component in reflection light from the work WK are received by a light position detection device such as a PSD via the light reception lens 104. When the work WK is displaced in the direction shown by the arrow X, the position of a light spot moves on the light reception surface of a light position detection device 105. Two output signals according to the position of a light spot on the light reception surface are output from the light position detection device 105 and are subjected to current-voltage conversion in current-voltage converting circuits (I-V converting circuits) 106a and 106b. One of the output signals of the optical position detection device 105 has a current value proportional to the distance from an end of the light reception surface to the light spot, and the other output signal has a current value proportional to the distance from the other end of the light reception surface to the light spot. Therefore, on the basis of the current values of the two output signals, the displacement of the work WK can be detected.\nIn such an optical displacement meter, the intensity of reflection light varies according to the material of the work WK and the surface state of the work WK. Consequently, the light output of the laser diode 102 has to be adjusted so that the light reception amount of the light position detection device 105 becomes a certain level. FIG. 94 is a block diagram showing an example of a conventional control circuit for controlling the light reception amount in the light position detection device 105. The control circuit of FIG. 94 includes the current-voltage converting circuits 106a and 106b, an adder 112, a subtracter 113, an error integration circuit 114, a reference voltage generating circuit 115, and a light output adjustment circuit 111. The current-voltage converting circuits 106a and 106b convert the current signal of the light position detection device 105 to a voltage signal. The adder 112 adds voltage signals on the far side and the near side, and outputs the light reception amount of the light position detection device 105 as a light reception amount voltage VL. The reference voltage generating circuit 115 generates a predetermined reference voltage Vr. The subtracter 113 outputs the difference between the light reception amount voltage VL obtained by the adder 112 and the reference voltage Vr generated by the reference voltage generating circuit 115. The error integration circuit 114 integrates error signals VE output from the subtracter 113, and provides the integrated error signals as a control voltage VC to the light output adjustment circuit 111. The voltage of the light output adjustment circuit 111 is controlled so that the light reception amount voltage VL output from the adder 112 becomes equal to the reference voltage Vr generated by the reference voltage generating circuit 115. Therefore, by controlling the voltage of the light output control signal Va provided from the light output adjustment circuit 111 to the drive circuit 101 in FIG. 93, the light reception amount in the light position detection device 105 can be controlled to a predetermined level.\nIn such an optical displacement meter, the reception light amount largely fluctuates and corresponds to the reflectance of light of a work according to the color, roughness, angle, and the like of the surface of a work. When the reception light signal is too small, or too large due to saturation or the like, the measurement accuracy deteriorates. Consequently, a technique of performing a feedback control that adjusts the light emission amount of a light emission device and the amplification factor (gain) of an amplifier so that the peak value of a light reception amount (image signal level) becomes a target value is developed (refer to, for example, Japanese Unexamined Patent Application Publication No. 2006-010361). An optical displacement meter of the technique has, as shown in FIG. 95, a light emitting device 102B for irradiating the work WK with light, an image sensor 105B for receiving light from the work WK and generating an image signal, a signal processing circuit including an amplifier 146B for amplifying the image signal from the image sensor 105B, and a controller 144B that executes feedback control of at least one of operation amounts including a light emission amount of the light emitting device 102B and the amplification factor of the amplifier 146B on the basis of the image signal from the signal processing circuit. The variable range of at least one of the operation amounts in the feedback control can be changed. The controller 144B sets a proper variable range of the operation amount on the basis of data of the operation amount in a predetermined period in a variable range setting mode. As a result, as shown in FIG. 96, while keeping the advantage of performing the feedback control of the light emission amount of the light emitting device 102B and the amplification factor of the amplifier 146B, the technique realizes higher speed of measurement.\nOn the other hand, an apparatus for measuring the shape (profile) of a work using the principle of light sectioning is developed. The light sectioning is obtained by two-dimensionally expanding the principle of triangulation as shown in FIGS. 97A and 97B. Specifically, as shown in FIG. 97A, a laser beam LB emitted from a sensor head SH to the work WK is set in a band shape, and the triangulation is increased in the width direction (X direction). As a result, the triangulation measuring technique is enlarged in the width direction as shown in FIG. 97B. Consequently, line-shaped (linear) light receiving devices are used in the triangulation measuring technique. On the other hand, light receiving devices JS arranged two-dimensionally are used in the light sectioning method.\nBy using an optical displacement meter for measuring a profile by using the principle of light sectioning, a stereoscopic shape of a work is measured. In the case of measuring a stereoscopic shape, there is a case such that intensity of reception light largely fluctuates according to the surface state and angle of a work. Conventionally, as the light receiving device, a CCD or CMOS sensor whose sensitivity characteristic (light receiving characteristic) is linear is used. Consequently, when the light amount changes considerably in a work, multiple exposure process is performed. In the process, while changing the amount of projected light and shutter time on the reception light side, imaging is performed on the work a plurality of times. From a plurality of reception light images obtained, ranges in which the light amount is proper are extracted. By combining the extracted images, a single image is obtained.\nIn the multiple exposure, however, the imaging process is performed a plurality of times. There is a problem such that the time required for obtaining one profile is long. For example, a reception light image requiring 15 to 16 ms per exposure is captured five to ten times successively, so that it takes about 0.5 to 1 second. Consequently, the technique is not suitable for an application requiring immediacy such as an inline process. When a work is position-deviated during image capturing a plurality of times, an error occurs at the time of composition of the multiple exposures. In such a case, an error can be compensated via software. In this case, the process amount further increases."}
{"text": "Food or beverage products that contain froth or foam are popular with consumers. For example, both hot and cold beverages, such as cappuccino, milkshakes, fruit-flavored drinks, and the like, may be prepared to have froth or foam. Recently, consumers increasingly prefer instant food and beverage products because of the convenience they provide. Such instant food and beverage products typically include a soluble powdered or granular food product that is dissolvable or dispensable in water or other fluid medium to form food or beverage products that may be consumed.\nRecent advances have provided methods and equipment for entrapping pressurized gases within internal voids of the powdered or granular food and beverage products. In this manner, upon dissolving or dispensing the food or beverage product in fluid medium, the pressurized gas is released forming a foam or froth. Thus, a frothy or foamy food or beverage can be prepared from the instant product.\nCurrent methods to prepare such powders may include a base product that has been pretreated to form a powder having internal voids. For example, the base product may be a carbohydrate powder that is prepared using gas-injected spray drying. The powder may be formed by spray drying a solution of the beverage product while injecting small gas bubbles into the solution stream to create a powder with internal voids.\nIn order to maximize the amount of foam that is created when the food or beverage powder is combined with fluid, the powder can be subjected to a pressurization process to entrap pressurized gas within the internal voids. The base powder substance is first placed into a batch vessel that is capable of withstanding relatively high temperature and pressure variations and the vessel is sealed to the external atmosphere. The vessel is next pressurized and heated by injecting pressurized inert gas into the vessel and applying heat to the walls of the vessel.\nAfter heating, the vessel is usually cooled by cooling the outer walls to slowly decrease the temperature of the powder from the heated temperature to a cooling temperature. The vessel is typically slowly cooled to ensure that it is not damaged due to a rapid change in temperature which could otherwise create thermal stresses in the vessel material. The temperature of the powder is usually held at the cooling temperature for an amount of time so that the powder particles become less permeable to the flow of gas. Once the powder is cooled, while still under pressure, the gas is sealed in the particle voids. Finally, the vessel is de-pressurized.\nWhile previous methods are generally sufficient for forming powdered and granular food and beverage products with entrapped pressurized gas, they do present some shortcomings. For example, the process time for a single batch of powder within a vessel in the prior methods is relatively high due, in part, to the lengthy time required to pressurize, heat, cool, and dc-pressurize the vessel. For the same reason, the energy requirements for each batch is high due to the significant amount of energy consumed each time the vessel is heated and subsequently cooled."}
{"text": "This invention relates to characterising the performance of a robot.\nA typical robot arm comprises a number of rigid limbs, each of which is connected to the next by an articulation which may, for example, be a revolute joint. A motor or other drive mechanism is provided to drive relative motion of the limbs at each joint.\nFIG. 1 shows an example of a drive mechanism at a robot arm articulation. The articulation 1 of FIG. 1 is located between a proximal arm member 2 and a distal arm member 3. The arm members are coupled together by a revolute joint defined by a rotation axis 4. An electric motor 5 is mounted in the proximal arm member 2. The output of the electric motor is a drive shaft 6 which connects to the input of a reduction gearbox 7. The gearbox 7 is also mounted in the proximal arm member. The output of the gearbox is a worm shaft 8. Worm shaft 8 engages a drive gear 9. Drive gear 9 is fast with the distal arm member 3 and is disposed about the rotation axis 4. This arrangement allows the motor 5 to drive relative motion of the arm members 2, 3 about the joint. In order to help control the motion of the arm joint, the joint is provided with a position encoder 11 and a torque sensor 10. The position encoder senses the configuration of the joint 1. The torque sensor senses the amount of torque applied about the joint. These values can be fed to a control processor and used as inputs to an algorithm to derive the electrical inputs to the motor 5.\nThe control processor may be required to control the joint to move into a desired configuration at a particular time. To achieve this the control processor needs to apply an appropriate amount of power to the motor. The power required depends on the inertia of the portion of the arm distal of the joint 1, the distance to be moved and the time over which the movement is to take place. The control processor may be capable of determining these in a known way. Additionally, however, the power required depends on the efficiency of the mechanical linkage between the motor 5 and the drive gear, involving the gearbox 7 and the transmission from the worm 8 to the drive gear 9. This mechanical efficiency is dependent on the meshing adjustment of the gearbox, the level of wear of the gearbox, which will vary over time, and the performance of any lubricant in the drivetrain. If the efficiency of the gearbox is not known then the arm can still be driven to the desired configuration using closed-loop control based on the output of position encoder 11. However, the arm may overshoot the desired configuration and have to be brought back, or fail to reach the desired position by the required time. Failures of the first type are especially important when the arm is part of a surgical robot because they may result in tissue of a patient being unnecessarily disrupted.\nAnother form of difficulty arises when the torque sensor is being used to sense the torque being applied on the joint 1 by an external force, such as gravity, at the same time as the distal limb is being driven by the motor 5. A torque will be measured about the joint. Part of that torque will be due to the external force and part will be due to the torque applied by the motor. Without knowing the performance of the gearbox 7 it is difficult to estimate the external torque. This may be problematic if the control strategy of the arm is dependent on an estimate of the externally applied torque.\nIt would be possible to measure the efficiency of the gearbox by providing an additional torque sensor at the input to the gearbox. However, this adds weight, cost and complexity to the arm.\nAnother alternative would be to remove the gearbox from the arm and measure its losses, but that is time-consuming and can only be undertaken when the arm is not in use.\nThere is a need for an improved method for estimating the efficiency of a mechanical linkage in a robot drivetrain."}
{"text": "hGH is a protein of 191 amino acids produced in the anterior lobe of the hypophysis throughout the life of an individual, and in greater quantities during the preadult period.\nThe growth hormone is synthesised in the form of a precursor and, once processed is secreted from the cell.\nhGH has for some time been used for curing some forms of dwarfism, which are due to lack of the hormone, and can also be used in the treatment of obesity and for healing burns and wounds.\nUntil a few years ago, the only source of the hormone was the hypophyses of corpses from which it was extracted at low yields by a complex and expensive method.\nMethods have recently been developed for the preparation of hGH by means of fermentation, with the use of host organisms transformed by recombinant DNA techniques.\nIn particular, GB patent no. 2055982 and EP patent no. 20147 described and claim a method for the preparation of hGH with the use of an engineered strain of Escherichia coli (E. coli) containing a plasmid which includes the structural sequence of the DNA which codes for the mature protein of 191 amino acids. The sequence is located downstream of a promoter and a ribosome recognition site (RBS) which are necessary for the transcription and translation processes, and juxtaposed to the ATG triplet which codes for methionine.\nThe presence of the methionine is necessary since it represents the starting signal for the translation of the whole protein.\nKnown methods, therefore produce a product which is constituted by the amino acid sequence of natural hGH with a methionine (Met) on its amino-terminal.\nAlthough the presence of this amino acid does not seem to affect the activity of the hormone (Olson K. C. et al. (1981), Nature 293, 408), immunological data show that in a fairly high percentage of patients treated with the Met:hGH the appearance of antibodies against the hormone, which prevent extended therapeutic treatment, can be shown.\nThere is consequently a need to provide a protein identical to the natural one.\nVarious methods have been proposed in the art for the removal in vivo and in vitro of methionine from the recombinant proteins, and in particular from hGH. On of the in vivo methods is based essentially on the use of signal sequences responsible for the activation of the transport of proteins through the cell membrane of Escherichia coli (E. coli) or Bacillus subtilis (B. subtilis).\nIn general this method comprises the construction of hybrid plasmids containing the nucleotide sequence which codes for mature hGH, fused at its 5' terminal to a sequence which codes for a signal sequence (a leader sequence) and the transformation of host cells by the hybrid plasmids. Suitably grown, the cells produce a fused protein, that is a protein constituted by hGH and the signal sequence. The sequence is then removed at the membrane level by a specific endopeptidase and the mature protein released into the periplasmic space or into the outside environment. In particular, the leader sequences of the protein OmpA (Hsiung H. M. et al. (1986) Biotechnology 4, 991) and the endotoxin of E. coli were used for the preparation of mature hGH from transformed E. coli cells and the signal sequence of the neutral protease of B. amyloliquefaciens was used for the preparation of hGH from B. subtilis cells.\nHowever, these known methods have disadvantages resulting on the one hand from the use of E. coli, which is a pathogenic organism in man, and on the other hand from the production of an hGH containing from 1 to 4 residual amino acids at its amino-terminus due to the partial processing of the protein expressed by B. subtilis.\nA second method for removal of the methionine in vitro consists of modification of the gene which codes for a protein so that its synthesis product is constituted by the natural protein and by a longer or shorter peptide sequence which can be removed by treatment with enzymes. In general, this method suffers from the fact that if residual amino acids, which are subject to attack by the enzymes used for the hydrolysis are present in the protein concerned, a greater or lesser percentage of the protein itself is degraded, consequently lowering the production yield.\nIt is therefore fundamental to provide a specific treatment which enables the fused protein to be hydrolysed exclusively at the point of connection between the product of interest and the amino-terminal peptide.\nIt is known that Factor Xa, a blood serine protease involved in the complex process of clotting, recognises the Glycine-Arginine (Gly-Arg) sequence, and in particular has a very great affinity for the isoleucine-glutamic acid-glycine-arginine (Ile-Glu-Gly-Arg) tetrapeptide.\nRecently, Nagai and collaborators (Nagai, K. et al (1985), Proc. Natl. Acad. Sci. USA 82,7252 have shown that fused products, in which the amino-terminal portion had been positioned to protect the protein of interest from the action of microbial proteases, were processed correctly and specifically when the Ile-Glu-Gly-Arg sequence was positioned at the connecting point.\nWhilst enabling the correct processing of the fused proteins and therefore the production of hGH in the natural form, these known methods however have disadvantages resulting from the difficulty of purification of the product obtained.\nIn fact, the precursor of hGH is difficult to separate both from the pool of proteins produced by the recombinant cells and from the hormone obtained after enzymatic hydrolysis."}
{"text": "The present invention relates to gas or combustion turbine apparatus, gas turbine electric power plants and control systems and operating methods therefor.\nIndustrial gas turbines may have varied cycle, structural and aerodynamic designs for a wide variety of uses. For example, gas turbines may employ the simple, regenerative, steam injection or combined cycle in driving an electric generator to produce electric power. Further, in these varied uses the gas turbine may have one or more shafts and many other rotor, casing, support and combustion system structural features which can vary relatively widely among differently designed units. They may be aviation jet engines adapted for industrial service as described for example in an ASME paper entitled \"The Pratt and Whitney Aircraft Jet Powered 121MW Electrical Peaking Unit,\" presented at the New York Meeting in November-December 1964.\nOther gas turbine uses include drive applications for pipeline or process industry compressors and surface transportation units. An additional application of gas turbines is that which involves recovery of turbine exhaust heat energy in other apparatus such as electric power or industrial boilers or other heat transfer apparatus. More generally, the gas turbine air flow path may form a part of an overall process system in which the gas turbine is used as an energy souroce in the flow path.\nGas turbine electric power plants are usable in base load, mid-range load and peak load power system applications. Combined cycle plants are normally usable for the base or mid-range applications while the power plant which employs a gas turbine only as a generator drive, typically is highly useful for peak load generation because of its relatively low investment cost. Although the heat rate for gas turbines is relatively high in relation to steam turbines, the investment savings for peak load application typically offsets the higher fuel cost factor. Another economic advantage for gas turbines is that power generation capacity can be added in relatively small blocks such as 25MW or 50MW, as needed, for expected system growth thereby avoiding excessive capital expenditure and excessive system reserve requirements. Further background on peaking generation can be obtained in articles such as \"Peaking Generation\" a Special Report of Electric Light and Power dated November 1966.\nStartup availability and low forced outage rates are particularly important for peak load power plant applications of gas turbines. Thus, reliable gas turbine startup and standby operations are particularly important for power system security and reliability.\nIn the operation of gas turbine apparatus and electric power plants, various kinds of controls have been employed. Relay-pneumatic type systems form a large part of the prior art, but have heretofore not provided the flexibility desired, particularly in terms of decision making. Furthermore, such prior art systems have been characterized by being specially designed for a given turbine plant, and accordingly are not adaptable to provide different optional features for the user. More recently, electronic controls of the analog type have been employed as perhaps represented by U.S. Pat. No. 3,520,133 entitled Gas Turbine Control System and issued on July 14, 1970 to A. Loft or by the control referred to in an article entitled Speedtronic Control, Protection and Sequential System and designated as GER-2461 in the General Electric Gas Turbine Reference Libary. See also U.S. Pat. No. 3,662,545, which discloses a particular type of analog acceleration control circuit for gas turbine; U.S. Pat. No. 3,340,883, relating to an analog acceleration, speed and load control system for a gas turbine. A wide variety of controls have been employed for aviation jet engines including electronic and computer controls as described for example in a March 1968 ASME Paper presented by J. E. Bayati and R. M. Frazzini and entitled Digatec (Digital Gas Turbine Engine Control), an April 1967 paper in the Journal of the Royal Aeronautical Society authored by E. S. Eccles and entiteld \"The Use of a Digital Computer for On-Line Control of a Jet Engine,\" or a July 1965 paper entitled \"The Electronic Control of Gas Turbine Engines\" by A. Sadler, S. Tweedy and P. J. Colburn in the July 1967 Journal of the Royal Aeronautical Society. However, the operational and control environment for jet engine operation differs considerably from that for industrial gas turbines.\nThe aforereferenced U.S. application W.E.40,062, assigned to the present assignee, presents an improved system and method for operating a gas turbine with a digital computer control system. In this system, one or more turbinegenerator plants are operated by a hybrid digital computer control system, wherein logic macro instructions are employed in programming the computer for logic operations of the control system.\nIn referencing prior art publications or patents as background herein, no representation is made that the cited subject matter is the best prior art.\nWhile industrial gas turbine apparatus gas turbine power plants have attained a great sophistication, there remain certain operational limitations in flexibility, response speed, accuracy and reliability. Further limits have been in the depth of operational control and in the efficiency or economy with which single or multiple units are placed under operational control and management. Limits have existed on the economics of industrial gas turbine application and in particular on how close industrial gas turbines can operate to the turbine design limits over various speed and/or load ranges.\nIn gas turbine power plants, operational shortcomings have existed with respect to plant availability and load control operations. Compressor surge control response has been limited, particularly during startup. Temperature limit control has been less protective and less responsive than otherwise desirable.\nGenerally, overall control loop arrangements and control system embodiments of such arrangements for industrial gas turbines have been less effective in operations control and systems protection than is desirable. Performance shortcomings have also persisted in the interfacing of control loop arrangements with sequencing controls.\nWith respect to industrial gas turbine startup, turbine operating life has been unnecessarily limited by conventional startup schemes. Sequencing systems have typically interacted with startup controls less effectively than desirable from the standpoint of turbine and power plant availability. More generally, sequencing systems have provided for systematic and protective advance of the industrial gas turbine operations through startup, run and shutdown, but in doing so have been less efficient and effective from a protection and performance standpoint than is desirable.\nRestrictions have been placed on operations and apparatus management particularly in gas turbine power plants in the areas of maintenance and plant information acquisition. Further management limits have existed with respect to plant interfacing with other power system points, operator panel functionality, and the ability to determine plant operations through control system calibration and parameter changes.\nThe computerized gas turbine control as disclosed in W.E. application Ser. No. 319,114 has been highly successful in providing control capability and flexibility of control options that had not previously been incorporated into an all hardware type system. However, while the computerized, or software control system provides substantial advantages due to its logic performing capability, historical data storage and diagnostic programs, it also has a number of shortcomings. The interface between the turbine and its associated analog signals and the computer controller presents areas for future development and improvement. The analog input system is a complex multiplexing arrangement requiring sharing of the scan time by the variables which must be scanned or read \"independently\". In the system disclosed, there is a scanning rate of 30 per second, meaning that 30 input variables per second can be read, imposing a limitation on the ability of the system to respond rapidly to a given input variable when program running time is also added to the delay. In addition, the computer system itself incorporates elaborate techniques of self-diagnosis of failure, which can result in turbine shut-downs when the computer has determined that something has failed within the central processor, input-output, or peripheral hardware. It is most difficult for the computer to determine whether the failure is of a sufficiently critical nature to require shutdown. In fact, it has been found that failures in the analog input-output system may not be readily differentiated, leaving the computer no choice but to shut down the entire turbine system for a failure which may not justify loss of load availability. Since all monitoring and protection paths are channeled through a central processor, a self-determination of failure in the central processor, analog input multiplexing or output system by the computer controller necessitates blocking off all channels, such that complete system shutdown is required. Furthermore, even during normal operation, the computerized system provides low visibility with respect to the health of the control system. The essential intermixing of the control paths through the central processor makes it difficult for the operator to obtain information as to the mode of control at any moment, or to obtain quantitative information as to the relative magnitudes of the different control signals. In short, the increased flexibility of the software system is achieved at the expense of operator visibility such as permits optimum maintenance procedures. Accordingly, there is a great need in the art for a turbine system having a control with the logic capability of a digital system, but retaining the advantages which are inherent in simpler designs.\nThe gas turbine control system as disclosed herein incorporates novel features which are specifically designed to meet the above general requirements, and which go further in providing operating capability not heretofore available in any turbine control system. The control system of this invention includes a plurality of continuously closed control loops, each of which continuously generates a control signal adapted to control the turbine fuel system, and thereby control available fuel to the turbine, thereby controlling turbine operation itself. Each of the control loops contains logic capability, is adaptable to be constructed in different hardware forms, and provides continuous visual indication for the operator and continuous monitoring for alarm or turbine shutdown. In this manner, should failure, or even a lesser malfunction, occur in any of the control paths, a backup control signal is available to take over turbine control, without the failure causing loss of turbine availability. Furthermore, means are provided by which the operator can immediately determine the source and, in many instances, the nature of the malfunction, so that corrective maintenance can be quickly and efficiently undertaken.\nAnother specific improvement is the provision of adapting the turbine control for changes in ambient temperature, such as occur between summer and winter operation. In prior art systems, which are dependent solely on monitoring of internal turbine conditions, unwanted operating limitations are imposed by changes in ambient temperature. Such limitations have been reduced substantially by the novel adaptive control means disclosed herein.\nAnother area of great importance in gas turbine contnrol is that of immediately meeting load demand upon generator breaker closing. Past controls have generally provided for a continuous buildup of load, starting from zero load at generator breaker closing and proceeding roughly linearly to a desired load level. However, there are a number of applications where it is required, or at least highly desirable, to provide an essentially instantaneous pickup of load. Accordingly, this invention provides novel means for controlling the turbine fuel flow so as to provide capability for such immediate load pickup.\nA critical portion of the operation of any turbine involves the starting sequence, at which time the turbine undergoes severe temperature changes, with possible resulting damage due to thermal stress. The turbine control system of this invention accordingly incorporates novel features to limit turbine speed change as a function of monitored turbine temperature, and to schedule the fuel supplied to the turbine combustor element so as to minimize risk of thermal damage during the starting operation. One of the novel techniques employed in this respect is the specific means by scheduling bypass fuel flow in the turbine fuel system during startup, so as to control the fuel pressure at the combustor nozzles. A bypass temperature limiter valve suitable in this operation, with which desired combustor nozzle fuel pressure is obtained during ignition, is disclosed in co-pending U.S. Pat. application Ser. No. 261,192, assigned to the same assignee. The technique disclosed herein involves novel means of utilizing apparatus such as is disclosed in Ser. No. 261,192.\nOne of the greatest needs in any turbine control system is that of providing operator flexibility, and in particular providng the operator with the capability of efficiently changing load as desired. Most prior art turbine control systems are quite limited in the degree of flexibility available to the operator, e.g., only discrete operating load levels are available, or the available means of changing the load level to a desired level is cumbersome and/or cannot be achieved at a desired rate. Accordingly, the control system of this invention provides novel means having essentially unlimited flexibiliity for operator change of the load level, rate of change of such load level, the ability to hold load at any desired level, the ability to return to any predetermined load level. This capability is constantly backed up by temperature control capability, such that no matter what the operator attempts to do, operational limits are automatically imposed as a function of sensed turbine temperatures.\nA yet further need in a turbine control system is that of providing a reliable and workable monitoring system. As pointed out hereinabove, in computerized turbine control systems all monitoring and protection paths are channeled through a central processor, which frequently results in the computer requiring complete shutdown when, in fact, the turbine is being operated within safe limits. Also, even during normal operation, the computerized system frequently does not permit the degree of monitoring visibility which is highly advantageous for providng the operator with optimum ability to oversee the turbine operation. In order to overcome these difficulties and provide improved visibility and reliability, while maintaining maximum turbine availability, there is a need for a control system designed so as to provide continuous visual indication as to the current mode of turbine control, so that the operator can determine the health of the control system. This may be achieved by providing discrete modularized control paths which are in constant communication with the turbine and which generate independent control signals, and means for determining and displaying which of said independent signals is at any given moment in control of the turbine operation.\nTurbine availability may also be greatly enhanced by providing means for automatically restarting the turbine after an automatic shutdown, and upon a determination that safe conditions exist for such restarting. In many cases, the condition which caused the turbine to be placed in shutdown is corrected, or corrects itself, shortly after shutdown is initiated. However, in conventional turbine control systems, the turbine must be brought substantially completely to the shutdown state, and then restarted all over again. It is clear that this results in avoidable loss of turbine availability, and that there is a need to minimize the loss of availability by restarting as soon as turbine conditions permit. The gas turbine monitoring system as disclosed herein provides a shutdown subsystem having a novel arrangement for automatically restarting the turbine after correction of the malfunction which caused shutdown.\nThe monitoring system of this invention also incorporates a unique system and method for optimizing load availability while providing an alarm or other means to alert the operator to the existence of a control malfunction which must be corrected but which does not merit immediate turbine shutdown. This unique system avoids the inflexibility of currently used monitoring circuits which are designed either to be fail-safe (in which case load availability is sacrificed to ensure shutdown), or which are designed to fail in a designated direction, (thus always proividing continuing availability, but at the cost of not shutting down in instances where shutdown might be required)."}
{"text": "The present invention relates to high temperature insulating, particularly in connection with providing a metallurgical soaking to a cooling ingot of steel to slow down the rate of cooling, enabling self-heat induced thermal structural homogenization within the ingot so that upon subsequent mill processing (rolling, swaging, forging or the like), equipment breakage is thereby minimized. This is accomplished through the present invention through structure affording the additional advantage of avoiding the high volume and energy usage of conventional soaking pits used for such purposes.\nThe insulating apparatus is usable for furnaces generally operating in a similar range of temperatures.\nDuring steel production, molten steel (at about 1500.degree. C.) is poured into iron molds and allowed to freeze on the outside surfaces. The mold is then stripped away and cooling is continued until the ingot is frozen. In order to permit hot rolling, the ingot is then put in a soaking pit which has burners therein (usually oil fired) and resides in such pit for half a day, exposed to flame or to air heated by the flame to come back up to appropriate temperature for rolling, in the range of 1000-1400.degree. C. The soaking pit structure is quite bulky and one or more ingots \"soaked\" therein generally occupy a small portion of the space. But the entire space has to be heated up and the energy consumption involved is substantial. The handling equipment associated with the soaking pit is also voluminous and expensive.\nIt is an important object of this invention to eliminate the energy consumption associated with such soaking process.\nIt is a further object of the invention to reduce the volume, mass and expense of the soaking equipment and associated handling equipment and process steps.\nIt is a further object of the invention to reduce scaling of soaked metals.\nIt is a further object of the invention to reduce initial freezing rate of the ingot with resultant improvement of ingot grain structure.\nIt is a further object of the invention to provide insulating method and apparatus which enables self-heat induced soaking of a poured ingot.\nIt is a further object of the invention to provide insulation apparatus and process usable broadly in connection with 900.degree. C. or above applications.\nIt is a further object of the invention to provide practical and effective high temperature insulation.\nThe foregoing objects also include the further ancillary object of achieving each such objects in combination with one or more of the others."}
{"text": "The present invention relates to an enzymatic process for producing cytidine diphosphate choline (hereinafter referred to as CDP-choline) which is useful as a medicine.\nCDP-choline is a biosynthetic intermediate for phosphatidyl choline (lecithin), which is a phospholipid, and is useful for the treatment of head injuries, disturbance of consciousness following cerebral surgery, Parkinson's disease, postapoplectic hemiplegia, etc.\nTwo kinds of processes are known for producing CDP-choline: chemical synthetic processes disclosed in Japanese Published Examined Patent Applications Nos. 6541/64, 1384/67, 6558/88, etc.; and enzymatic processes utilizing cells of microorganisms such as yeast disclosed in Japanese Published Examined Patent Applications Nos. 2358/73, 40757/73 and 40758/73 and Japanese Published Unexamined Patent Applications Nos. 109996/78, 14593/79 and 313594/88.\nA feature common to these known processes is the use of cytosine nucleotides such as cytidine-5′-monophosphate (hereinafter referred to as CMP), cytidine-5′-diphosphate (hereinafter referred to as CDP), cytidine-5′-triphosphate (hereinafter referred to as CTP), and cytosine, or their precursors as starting materials.\nOf these starting materials, the basic starting material, CMP, is mainly produced by the RNA (ribonucleic acid) decomposition method which provides four types of nucleotides at the same time. This method is inefficient in that it is impossible to selectively obtain CMP."}
{"text": "An integrated circuit is typically formed on a substrate by the sequential deposition of conductive, semiconductive, or insulative layers on a silicon wafer. A variety of fabrication processes require planarization of a layer on the substrate. For example, for certain applications, e.g., polishing of a metal layer to form vias, plugs, and lines in the trenches of a patterned layer, an overlying layer is planarized until the top surface of a patterned layer is exposed. In other applications, e.g., planarization of a dielectric layer for photolithography, an overlying layer is polished until a desired thickness remains over the underlying layer.\nChemical mechanical polishing (CMP) is one accepted method of planarization. This planarization method typically requires that the substrate be mounted on a carrier or polishing head. The exposed surface of the substrate is typically placed against a rotating polishing pad. The carrier head provides a controllable load on the substrate to push it against the polishing pad. A polishing liquid, such as abrasive slurry, is typically supplied to the surface of the polishing pad.\nOne problem in CMP is determining whether the polishing process is complete, i.e., whether a substrate layer has been planarized to a desired flatness or thickness, or when a desired amount of material has been removed. Variations in the slurry distribution, the polishing pad condition, the relative speed between the polishing pad and the substrate, and the load on the substrate can cause variations in the material removal rate. These variations, as well as variations in the initial thickness of the substrate layer, cause variations in the time needed to reach the polishing endpoint. Therefore, determining the polishing endpoint merely as a function of polishing time can lead to within-wafer non-uniformity (WIWNU) and wafer-to-wafer non-uniformity (WTWNU).\nIn some systems, a substrate is optically monitored in-situ during polishing, e.g., through a window in the polishing pad. However, existing optical monitoring techniques may not satisfy increasing demands of semiconductor device manufacturers."}
{"text": "This invention relates to new and useful improvements in automatically aligning hitch assemblies designed to ease the connection of a towed device to a towing vehicle and although the present specification and drawings relate primarily to the connection between a tractor draw bar and an implement hitch, nevertheless it is readily adapted for use between a truck and trailer or a towing vehicle to anything requiring to be towed.\nDifficulty exists in backing up the towing vehicle such as a tractor, to a towed device such as an implement, in a position so that the draw bar of the tractor, which is provided with an aperture, lines up exactly with the aperture in the corresponding clevis or hitch or the implement thus allowing a pin to be dropped through the tow bar and the hitch of the implement thereby connecting the two together. This problem is increased with the advent of tractor cabs being installed on modern tractors which cuts down the visibility.\nIt will be appreciated that the larger the equipment, the harder it is to judge due to the decreased visibility and distance from the connecting points. The same problem exists when attempting to hitch a truck or car to a trailer or implement as it is very difficult for the operator because the hitch is located at the centre of the rear bumper of the car or truck and is therefor not visible from the driver's seat.\nAttempts have been made to overcome this disadvantage and one such solution is the provision of a mirror on the towing vehicle adjacent the hitch through which the operator may view the hitch, but this is also unsatisfactory.\nVarious other solutions have been attempted and U.S. Pat. No. 2,791,443 discloses one such solution which utilizes toggle arms and collapsible guide means.\nU.S. Pat. No. 2,386,195 discloses a vehicle coupler which includes link means that collapse when the vehicle is coupled and U.S. Pat. No. 2,341,528 shows a similar structure.\nU.S. Pat. No. 1,575,773 discloses a tractor hitch which includes pivoted arms and spring means for connecting the hitch, but they all suffer from several disadvantages and do not provide the flexibility of the present assembly nor do they provide the automatic locking of the hitch assembly in the retracted position once the initial manual connection has been made."}
{"text": "The rapid increase in the incidence of gram-positive injections—including those caused by antibiotic resistant bacteria—has sparked renewed interest in the development of novel classes of antibiotics. One such class is the lipopeptide antibiotics, which includes daptomycin. Daptomycin has potent bactericidal activity in vitro against clinically relevant gram-positive bacteria that cause serious and life-threatening diseases. clinically relevant gram-positive bacteria that cause serious and life-threatening diseases. These bacteria include resistant pathogens, such as vancomycin-resistant enterococci (VRE), methicillin-resistant Staphylococcus aureus (MRSA), glycopeptide intermediary susceptible Staphylococcus aureus (GISA), coagulase-negative staphylococci (CNS), and penicillin-resistant Streptococcus pneumoniae (PRSP), for which there are very few therapeutic alternatives. See, e.g., Tally et al., 1999, Exp. Opin. Invest. Drugs 8: 1223-1238, hereafter “Tally”. Daptomycin's inhibitory effect is a rapid, concentration-dependent bactericidal effect in vitro and in vivo, and a relatively prolonged concentration-dependent post-antibiotic effect in vivo.\nDaptomycin is described by Baltz in Biotechnology of Antibiotics, 2nd Ed., ed, W. R. Strohl (New York: Marcel Dekker, Inc.), 1997. pp. 415-435, hereafter “Baltz.” Daptomycin, also known as LY 146032, is a cyclic lipopeptide antibiotic that can be derived from the fermentation of Streptomyces roseosporus. Daptomycin is a member of the factor A-21978C0 type antibiotics of S. roseosporus and is comprised of a decanoyl side chain linked to the N-terminal tryptophan of a cyclic 13-amino acid peptide (FIG. 1). Daptomycin has an excellent profile of activity because it is highly effective against most gram-positive bacteria: it is highly bactericidal and fast-acting; it has a low resistance rate and is effective against antibiotic-resistant organisms. The compound is currently being developed in a variety of formulations to treat serious infections caused by bacteria, including, but not limited to, methicillin resistant Staphylococcus aureus (MRSA) and vancomycin resistant enterococci (VRE).\nA number of United States Patents describe A-21978C-antibiotics and derivatives thereof including daptomycin (LY 146032) as well as methods of producing and isolating the A-21978C antibiotics and derivatives thereof.\nU.S. Pat. Re. 32333, Re. 32,455 and U.S. Pat. No. 4,800,157 describe a method of synthesizing daptomycin by cultivating Streptomyces roseosporus NRL15998 under submerged aerobic fermentation conditions. U.S. Pat. No. 4,885,243 describes an improved method of synthesizing daptomycin by feeding a fermentation culture a decanoic fatty acid or ester or salt thereof.\nU.S. Pat. Re. 32,310, Re. 32,311, U.S. Pat. Nos. 4,537,717, 4,482,487 and 4,524,135 describe methods of deacylating the A-21978C antibiotic and reacylating the peptide nucleus and antibiotic derivatives made by this process. All of these patents describe a purified deacylated A-21978C antibiotic nucleus or a derivative thereof which was isolated from the fermentation broth by filtration and then purified by Diaion HP-20 chromatography and silica gel/C18 chromatography.\nU.S. Pat. Re. 32,333 and Re. 32,455 disclose a purification method in which a filtrate of whole fermentation broth was purified through a number of precipitation and extraction steps to obtain a crude A-21978C complex. The crude complex was further purified by ion exchange chromatography on IRA-68 and two rounds of silica gel chromatography. Individual A-21978C factors were separated by reverse-phase silica gel or silica gel/CIS. U.S. Pat. Re. 32,333 and Re. 32,455 also disclose that A-21978C may be purified by batch chromatography using Diaion HP-20 resin followed by silica-gel column chromatography.\nU.S. Pat. No. 4,874,843 describes a daptomycin purification method in which the fermentation broth was filtered and passed through a column containing HP-20 resin. Alter elution, the semipurified daptomycin was passed through a column containing HP-20ss, and then separated again on HP-20 resin. The '843 patent, states that final resolution and separation of daptomycin from structurally similar compounds by this method is impeded by the presence of impurities that are not identifiable by ultraviolet analysis of the fermentation broth. The '843 patent further states that attempts to remove these impurities by reverse phase chromatography over silica gel, normal phase chromatography over silica gel or ion exchange chromatography also failed to significantly improve the purity of daptomycin. The '843 patent also discloses a “reverse method” for purification, comprising the steps of contacting an aqueous solution of the fermentation product with a non-functional resin in aqueous phase, physically removing the water from the charged resin, rewetting the charged resin with a polar organic solvent, washing the resin with the organic solvent, eluting the fermentation product from the resin by increasing the polarity of the solvent and recovering the fermentation product. The '843 patent teaches that this method improves the final purity from about 80% to about 93% and increases the yield from about 5% to about 35%; however, the '843 patent does not disclose the type of impurities present in the daptomycin preparation.\nU.S. Pat. No. 5,912,226 describes the identification and isolation of two impurities produced during the manufacture of daptomycin. Daptomycin, an α-aspartyl peptide, becomes transpeptidated to form a stable intermediate in which the aspartyl group becomes an anhydro-succinimido group (FIG. 3). The '226 patent teaches that the presence of this intermediate, designated anhydro-daptomycin, is more pronounced at pH 4-6. Rehydration of the anhydro-succmimido form produces a second degradation product that contains an β-aspartyl group and is designated the β-isomer form of daptomycin (FIG. 2).\nThe '226 patent discloses that the t-BOC derivative of anhydro-daptomycin may be isolated by chromatography over reverse phase silica gel/C-18 column, precipitated, and repurified by reverse phase silica gel/C-18 chromatography. The '226 patent also teaches that the β-isomer form of daptomycin may be purified by chromatography over a Diaion HP-20ss resin, desalted by chromatography over a Diaion HP-20 resin, and further purified using a reverse-phase C-18 column followed by a HP-20 resin column in reverse mode.\nKirsch et al. (Pharmaceutical Research, 6:387-393, 1989, hereafter “Kirsch”) stated that anhydro-daptomycin and the β-isomer were produced in the purification of daptomycin. Kirsch described methods to minimize the levels of anhydro-daptomycin and the β-isomer through manipulation of pH conditions and temperature conditions. However, Kirsch was unable to stabilize daptomycin and prevent the conversion of daptomycin to anhydro-daptomycin and its subsequent isomerization to β-isomer. Kirsch was also unable to prevent the degradation of daptomycin into other degradation products unrelated to anhydro-daptomycin and β-isomer.\nThe '226 patent states that daptomycin may be prepared using these procedures so that the daptomycin contains no more than 2.5% by weight of a combined total of anhydro-daptomycin and β-isomer, but gives no indication of the levels of other impurities. In the method taught in U.S. Pat. No. 4,874,843 and in large-scale preparations of daptomycin for clinical trials, the highest daptomycin purity levels observed has been about 90%-93%. There is a need for a commercially feasible method to produce more highly purified daptomycin and, if possible, to increase its yield after purification. Furthermore, it would be desirable to obtain purified daptomycin that contains little or none of anhydro-daptomycin and the β-isomer form of daptomycin. It would also be desirable to reduce the levels of a number of other impurities in daptomycin. However, there has been no method available in the art that has been shown to be able to further reduce the levels of anhydro-daptomycin, β-isomer form and other impurities in the daptomycin product."}
{"text": "Distribution of digital terrestrial television (DTT) and mobile digital television (MDTV) frequently utilizes single frequency networks (SFN). In an SFN, several transmitters simultaneously send the same signal over the same frequency channel.\nThe transmitters in an SFN must be synchronized to send their signals at the same time to avoid interference at the receiving antennas. This is commonly achieved by installing global positioning system (GPS) receivers at all transmitter sites. GPS receivers, however, may be easily intentionally or unintentionally jammed, or fail for other reasons such as equipment failure, and represent an additional cost in the network in terms of equipment and supervision. Further, the military control of the GPS may be an issue.\nAlso known are techniques for time synchronization of network nodes without utilization of GPS. For instance, the network time protocol (NTP) may be used to synchronize the clocks of network nodes to a master node or a reference clock using time stamps. However, the accuracy of NTP, at least in non-dedicated networks, is far too limited for the purpose of time synchronization in digital television (DTV) distribution networks.\nWO 2008/103170 A discloses a method of a network client for extracting a reference frequency carried in the physical layer of a network signal originating from a server, and for using it to stabilize an oscillator of a clock of the client. The method also includes determining a clock correction value based on a server time stamp and a client time stamp.\nU.S. Pat. No. 7,535,931 B discloses a two-way time transfer protocol for estimating a time error between the clocks of two network nodes."}
{"text": "The majority of golf balls commercially available today are of a solid construction. Solid golf balls include one-piece, two-piece, and multi-layer golf balls. One-piece golf balls are inexpensive and easy to construct, but have limited playing characteristics and their use is, at best, confined to the driving range. Two-piece golf balls are generally constructed with a solid polybutadiene core and a cover and are typically the most popular with recreational golfers because they are very durable and provide good distance. These golf balls are also relatively inexpensive and easy to manufacture, but are regarded by top players as having limited playing characteristics. Multi-layer golf balls are comprised of a solid core and a cover, either of which may be formed of one or more layers. These balls are regarded as having an extended range of playing characteristics, but are more expensive and difficult to manufacture than are one- and two-piece golf balls.\nWound golf balls, which typically included a fluid-filled center surrounded by a layer of tensioned elastomeric material and a cover, were preferred for their spin and “feel” characteristics but were more difficult and expensive to manufacture than solid golf balls. Manufacturers are continuously striving to produce a solid ball that concurrently includes the beneficial characteristics of a wound ball.\nGolf ball playing characteristics, such as compression, velocity, and spin can be adjusted and optimized by manufacturers to suit players having a wide variety of playing abilities. For example, manufacturers can alter any or all of these properties by changing the materials and/or the physical construction of each or all of the various golf ball components (i.e., centers, cores, intermediate layers, and covers). Finding the right combination of core and layer materials and the ideal ball construction to produce a golf ball suited for a predetermined set of performance criteria is a challenging task.\nEfforts to construct a multi-layer golf ball have generally focused on the use of one or two cover layers formed of ionomeric and/or polyurethane compositions. It is desirable, therefore, to construct a golf ball formed of a urethane or urea outer cover layer, at least two interior cover layers, and at least one core layer, according to the present invention. In particular, it is desired that this construction include a stiff, thermoplastic polyurethane or polyurea intermediate cover layer in conjunction with a stiff, resilient thermoplastic inner cover layer and thermosetting castable outer cover layer."}
{"text": "1. Field of the Invention\nThe present invention relates to downloading system and, more particularly, to a downloading system which can be used by, for example, many and unspecified users to download various sorts of information according to their needs.\n2. Description of the Related Art\nVarious sorts of software suitable for use in information apparatuses, e.g., application software for use in personal computers and contents of media (magazines, newspapers, music software, video software, game software, etc.) reproducible and/or editable by a personal computer, an audio and visual apparatus or the like, are now being provided for pay or free.\nHowever, a user cannot obtain one of such various sorts of software except by purchasing it in a software package in the form of a CR-ROM or the like or by connecting his or her personal computer to a network such as the Internet and downloading it.\nTherefore, there is a demand for a software acquisition system which can be easily used by a number of persons to obtain information according to their needs."}
{"text": "The presently described advancements relate to a multi-core processor interface function and topology, and, more particularly, to a multi-core processor interface leveraging a generated topology for providing decreased energy utilization and increased operational speed, and associated methodology.\nThe “background” description provided herein is for the purpose of generally presenting the context of the inventive field. Work of the presently named inventors, to the extent described in this section, as well as aspects of the background description which may not otherwise qualify as prior art at the time of filing, are neither expressly or impliedly admitted as prior art, but instead, provided in the spirit of teaching the overall context of the advancements claimed herein.\nAs software increases in complexity and resource consumption, a corresponding demand on hardware platforms is created. In answering this demand, single-core processors and the performance imperative of Moore's Law, have limited the ability to keep pace by corresponding increases in processing power (i.e., via increasing clock speeds). To overcome this problem, multi-processor architectures have been developed such that two or more processors are utilized on a same board to achieve performance gains. These multi-processor architectures are known as “multi-core”. Multi-core processors increase processing performance while reducing processor footprints and providing better power consumption and thermal characteristics.\nYet, most legacy application software is written for a single processor. Such application software will not automatically run faster on an embedded multi-core computer; in fact it may not run at all. This may be the case even if the software is multi-threaded, as in a single processor the threads will run sequentially whereas in a multi-core computer they may run concurrently, exposing synchronization problems not discovered in the single processor case. Programming multithreaded code often requires complex coordination of threads and can easily introduce bugs that are difficult to find due to the interleaving of processing on data shared between threads. To further complicate matters, the threads are ideally distributed among the multiple cores in a manner that loads the cores equally, to achieve maximum benefit (performance and power consumption).\nGiven the increasing emphasis on multi-core chip design, the extent to which software can be multithreaded to take advantage of these new chips is likely to be the single greatest constraint on computer performance in the foreseeable future.\nPresently, there is a need for a portable API, devoid of the aforementioned deficiencies, providing an optimized communications topology for multi-core processor applications."}
{"text": "The present invention is directed toward an improved supporting stand for supporting items of electronic equipment. Such a stand is frequently used to support extremely expensive or delicate electronic units, and therefore should be as strong and sturdy as is practicable. However, it should preferably also be capable of being stored in a minimum amount of space when not actually required for use. As will be explained hereinafter, the above requirements are not fully met by electronic equipment of conventional type, which are subject to vibration when sideways acting forces are applied to them, and which cannot be folded up in such a way as to minimize the storage space necessary for them. The present invention is designed to overcome these disadvantages of prior art stands for electronic equipment units."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid-discharge control method and a liquid discharging apparatus. More particularly, the invention relates to a liquid-discharge control method when discharging a liquid using a head according to a liquid discharging method utilizing the generation of a bubble by heat, and to a liquid discharging apparatus. The invention also relates to an ink-discharge control method when performing recording on a recording medium using ink as a liquid, and to an ink-jet recording apparatus.\n2. Description of the Related Art\nFrom among conventionally known ink-jet recording methods, a recording method (a bubble-jet recording method) has been widely known in which a bubble is generated by heating ink within one of ink discharging ports and a liquid channel communicating therewith (these two components will be hereinafter termed a xe2x80x9cnozzlexe2x80x9d) by heating means, such as a heater or the like, and a fine ink droplet is discharged from the ink discharging port onto a recording medium by the function of the bubble, in order to form an image by consecutively discharging ink droplets from corresponding ink discharging ports in the same manner. This method has been applied to printers, copiers and the like.\nRecording heads which adopt the bubble-jet recording method are suitable for high-speed recording and high-quality image recording, because it is easy to increase the number of nozzles mounted in the head (provision of multiple nozzles) and to provide high-density nozzles. Particularly, attempts to increase the recording speed of a printer, a copier or the like by providing multiple nozzles or increasing the driving frequency for the recording head are actively being tried. There are also attempts to realize high resolution by reducing the amount of droplets (the amount of discharge) of ink, and thereby to improve the quality of a recorded image to the high-quality level of photography.\nIn the bubble-jet recording, a recording method has been known in which, when performing recording by discharging ink from a discharging port by a bubble generated by heating the ink, the bubble is caused to communicate with the atmospheric air if the internal pressure of the bubble is negative (U.S. Pat. No. 5,218,376). According to this method, it is possible to prevent the generation of ink mist during ink splash or discharge, so that the recording medium or the inside of the apparatus is not stained with the ink. In addition, since ink between the generated bubble and the discharging port can be substantially entirely discharged and the amount of the discharged ink is determined by the shape of the nozzle and the position of the heater, it is possible to perform stable recording in which the amount of the discharged ink droplet is always constant.\nIn the conventional bubble-jet recording method, when the duty ratio of an image is high or the ambient temperature is high, the temperature of the recording head is raised, so that the discharging direction may slightly change.\nFIGS. 11A through 13F are schematic diagrams illustrating manners in each of which an ink droplet is discharged from an ink discharging port of a recording head. The recording head adopts a recording method in which a bubble generated during recording is caused to communicate with the atmospheric air.\nFIGS. 11A-11C are enlarged views of a nozzle of the recording head.\nIn FIGS. 11A-11C, there are shown an electrothermal transducer 31 for heating ink, an ink discharging port 32, an ink supply port 33, and a discharging-port plate 35. FIG. 11A is a side cross-sectional view of the ink discharging nozzle. FIG. 11B is a side cross-sectional view of the ink discharging nozzle shown in FIG. 11A as seen from a direction rotated by 90 degrees from the state shown in FIG. 11A. FIG. 11C is a top plan view of the ink discharging nozzle. Line 11A-11A xe2x80x2 shown in FIG. 1C is a line of cutting plane for providing the side cross section shown in FIG. 11A. Line 11B-11Bxe2x80x2 shown in FIG. 11C is a line of cutting plane for providing the side cross sectional view shown in FIG. 11B.\nFIGS. 12A-12C and 12D-12F are diagrams, each illustrating how ink is discharged in the ink discharging state of the present invention in which a bubble generated from the recording head discharges ink by communicating with the atmospheric air in a state of negative pressure, as seen from the direction shown in FIG. 11A.\nFIGS. 12A-12C illustrate how the ink is discharged when the temperature of the recording head is high. FIGS. 12D-12F illustrate how the ink is discharged when the same energy (driving signal) as in the case shown in FIGS. 12A-12C is applied to the electrothermal transducer in a state in which the temperature of the recording head is close to the room temperature.\nAs can be understood from FIGS. 12A-12F, the generated bubble is larger when the temperature is high than when the temperature is close to the room temperature, and the direction of ink discharge slightly changes depending on the difference between the temperatures of the recording head. It is considered that this is because the amount of ink slightly remaining in the vicinity of the ink discharging port differs depending on the size of the generated bubble, and a portion where the rear end of the ink leaves the ink discharging port differs depending on the difference in the amount of ink remaining in the vicinity of the ink discharging port.\nThat is, in FIG. 12B, consider the amounts of ink remaining at portions n and m, each surrounded by a circle. When the temperature of the recording head is high and therefore the size of the generated bubble is large, the amount of ink remaining at portion m is small, and as shown in FIG. 12C, the direction of the ink droplet discharged from ink remaining at portion n opposite to portion m of the ink discharging port slightly deviates from the center of the ink discharging port (indicated by a broken line in FIG. 12C). On the other hand, when the temperature of the recording head is low and therefore the size of the generated bubble is small, although a bubble tends to be generated slightly toward the ink supply port, this tendency is small, and, as shown in FIG. 12F, ink is discharged substantially rectilinearly along the center line of the ink discharging port.\nIn the case of the conventional recording head in which a bubble does not communicate with the atmospheric air, when the discharged ink droplet is separated from ink within the nozzle, the influence of the separated ink remaining in the vicinity of the ink discharging port does not cause any particular problem, because the ink returns to the inside of the nozzle. However, in the recording head having the above-described configuration in which a bubble generated by driving the electrothermal transducer is caused to communicate with the atmospheric air, the above-described phenomenon that, when the temperature of the recording head is high, the portion where an ink droplet is separated during ink discharge is in the vicinity of the inner wall of the ink discharging port and the discharging direction deviates, as shown in FIG. 12B, is observed.\nFIGS. 13A-13C, and 13D-13F are diagrams, each illustrating how ink is discharged from the same viewpoint as shown in FIG. 11B, in order to illustrate the state of ink discharge shown in FIGS. 12A-12F in further detail.\nAs in the case of FIGS. 12A-12C, FIGS. 13A-13C illustrate how the ink is discharged when the temperature of the recording head is high. FIGS. 13D-13F illustrate how the ink is discharged when the same energy (driving signal) as in the case shown in FIGS. 13A-13C is applied to the electrothermal transducer in a state in which the temperature of the recording head is close to the room temperature.\nAs can be understood from FIGS. 13A-13C, when the temperature of the recording head is high, the generated bubble is large, and asymmetrical with respect to the center of the ink discharging port. Accordingly, as shown in FIG. 13A, a difference tends to occur between the velocity vectors P and Q of the ink due to the growth of the bubble, and a difference in the amount of ink remaining in the inner wall of the ink discharging port tends to occur between portions S and T. As a result, deviation in the discharging direction occurs.\nOn the other hand, when the temperature of the recording head is close to the room temperature, the generated bubble is small, and relatively symmetrical with respect to the center of the ink discharging port. Accordingly, as shown in FIG. 13D, a difference is hardly produced between the velocity vectors P and Q of the ink due to the growth of the bubble, and as shown in FIGS. 13E and 13F, the ink discharging direction is rectilinear along the center line of the ink discharging port (as indicated by a broken line shown in FIG. 13F) on an average.\nAs studied above, in the conventional head in which a bubble communicates with the atmospheric air in a state of negative pressure, the ink discharging direction tends to deviate when the temperature of the head is raised. As a result, the position where the ink droplet adheres on a recording sheet deviates, thereby causing problems to be solved such that white stripes are produced on the recorded image, or unevenness in density is produced due to overlap of recorded dots.\nAn undisclosed technique relating to the present invention will now be described.\nAs a result of further detailed studies about the relationship between the manner of ink discharge and input energy in a recording head in which a bubble communicates with the atmospheric air during ink discharge, the phenomenon that a part of the discharged ink droplet falls onto the electrothermal transducer (see FIG. 14C) is observed. Recent detailed studies have shown that this phenomenon of falling of the ink droplet during ink discharge occurs when the power of bubble generation is smaller than the power of bubble generation during ink discharge in which the bubble communicates with the atmospheric air, as in the above-described cases shown in FIGS. 12D-12F and FIGS. 13D-13F, at the same head temperature.\nThat is, according to the recent knowledge of the assignee of the present application, even in an ink discharge method in which ink is discharged by causing the bubble to communicate with the atmospheric air when the inner pressure of the bubble is negative, when input energy to the electrothermal transducer is high and the power of bubble generation is high, as can be understood from FIGS. 12D-12F and FIGS. 13D-13F, a discharging state in which the liquid droplet is disconnected from the liquid remaining in the vicinity of the discharging port at a position near the end portion of the discharging port is provided. On the other hand, when input energy to the electrothermal transducer is low and the power of bubble generation is relatively low, a discharging state in which the liquid droplet is disconnected from the liquid falling onto the electrothermal transducer in the vicinity of the center line of the discharging port occurs, as shown in FIGS. 14A-14C.\nThe term xe2x80x9cfallxe2x80x9d of the liquid or ink used in this specification indicates not only dropping of the liquid or ink in the direction of gravity, but also the movement and adherence of a part of the liquid or ink to be discharged, onto the surface of the substrate where the electrothermal transducer is provided, irrespective of the direction of the head.\nEven in the discharging method having such a phenomenon of fall, the discharged ink droplet has an appropriate discharging speed, for example, 10-25 m/sec, which is sufficient as the discharging speed. Hence, the possibility that an ink droplet having an extremely high discharging speed is produced as in the case of ink discharge in which the phenomenon of fall does not occur is not present, and rebound of an ink droplet from the recording medium (the generation of mist) hardly occurs.\nFurthermore, since the disconnection of the liquid does not occur in the vicinity of the end portion of the discharging port, deviation of the discharging direction of the droplet, which occurs at random when the phenomenon of fall is not present, hardly occurs.\nEven if the phenomenon of fall occurs, this phenomenon does not influence the ink discharging speed. Hence, in ink discharge in which the bubble communicates with the atmospheric air, it is desirable to perform ink discharge by stably producing the phenomenon of fall, from the viewpoint of stabilization of the discharging direction of the droplet and suppression of the generation of mist.\nHowever, even when adjusting the energy (the driving signal) supplied to the electrothermal transducer so that the ink droplet falls during ink discharge at the room temperature, the phenomenon of fall may not occur, in some cases, as shown in FIGS. 13D-13F because the energy of bubble generation increases although the same energy is supplied to the electrothermal transducer, if the head temperature is raised. In such a case, the position where the droplet is disconnected during ink discharge differs as described above, so that the ink discharging direction is not constant. As a result, the position where the ink droplet adheres to the recording medium deviates, thereby causing the generation of white stripes on the recorded image or the generation of unevenness in density due to differences in the overlapped positions of recorded dots. These phenomena may cause degradation in the quality of the recorded image.\nFIGS. 15A and 15B are enlarged diagrams, each illustrating a longitudinal line recorded by a recording head having two columns of ink discharging nozzles. When recording a desired image by the recording head having such a configuration, since two columns of nozzles are present, recording is usually controlled so that ink is discharged by providing a time difference between the two columns of nozzles. In order to facilitate the following description, one column of nozzles is termed a column A, and another column of nozzles is termed a column B.\nIn each of the nozzle columns A and B, nozzles are arranged with a density of 300 dots per inch (300 dpi), and the columns A and B are arranged in a state of being shifted by {fraction (1/600)} inch in the direction of the columns with each other. The interval between the columns A and B is {fraction (1/600)} inchxc3x975 (for 5 pixels of 600 dpi). When driving such a recording head with a driving frequency of 10 kHz, a time difference for recording 5 pixels of 600 dpi may be provided between the discharging timings for the nozzle columns A and B. This value equals 100xc3x975 xcexcsec. Thus, as shown in FIG. 15A, when a one-dot longitudinal line is recorded, respective dots are arranged along the longitudinal line.\nFIG. 15B illustrates the positions of discharged dots when the discharging direction has slightly changed due to a change in the temperature of the recording head. Since the nozzle columns A and B are in a symmetrical positional relationship, a change in the ink discharging direction is enhanced. When the ink discharging direction has changed in the above-described manner, the recorded image is disturbed. Particularly when performing high-resolution recording, it is necessary to control the ink discharging position with higher accuracy.\nWhen the positions of adherence of discharged dots on the recording medium deviate as a result of changes in the ink discharging direction as shown in FIG. 15B, a texture may be produced in a pattern representing halftone. Usually, an error diffusion method, a dither method or the like is used as the method for representing pseudo-halftone. It is considered that the above-described texture is produced by synchronism of a specific pattern used in the above-described method with deviation in the positions of recorded dots. The generation of such texture also causes degradation in the picture quality.\nThe present invention has been made in consideration of the above-described conventional approach and undisclosed technique.\nIt is an object of the present invention to provide a recording control method and a recording apparatus in which very accurate recording positions can always be obtained irrespective of changes in the temperature of a recording head.\nAccording to one aspect of the present invention, a liquid-discharge control method, in which a bubble is generated by supplying a liquid with thermal energy by driving one of electrothermal transducers, and the liquid is discharged from a corresponding one of nozzles by causing the bubble to communicate with the atmospheric air, includes a detection step of detecting a temperature of a head where the electrothermal transducers and the nozzles are provided, an adjusting step of adjusting a width of a driving signal for driving the electrothermal transducer so that a part of the discharged liquid falls onto a side where the electrothermal transducer is provided, based on the temperature detected in the detection step, and a recording step of discharging the liquid by driving the head using the driving signal having the adjusted width adjusted in the adjusting step.\nAccording to another aspect of the present invention, a liquid discharging apparatus including a head, which includes electrothermal transducers, each for generating a bubble by supplying a liquid with thermal energy, and nozzles for discharging the liquid, for causing the bubble to discharge the liquid by causing the generated bubble to communicate with the atmospheric air, includes detection means for detecting a temperature of the head, adjustment means for adjusting a width of a driving signal for driving one of the electrothermal transducers so that a part of the discharged liquid falls onto a side where the electrothermal transducer is provided, based on the temperature detected by the detection means, and recording means for discharging the liquid by driving the head using the driving signal having the adjusted width adjusted by the adjustment means.\nThe foregoing and other objects, advantages and features of the present invention will become more apparent from the following detailed description of the preferred embodiments taken in conjunction with the accompanying drawings."}
{"text": "Wideband communication systems using optical signals now commonly employ wavelength division multiplexing (WDM) or dense wavelength division multiplexing (DWDM) to transfer and transmit vast amounts of data at very high data rates. In these multiplexing systems, multiple signal-bearing channels are centered on specified spaced apart wavelengths and sources (e.g. lasers) are set to those wavelengths. Whether combining signals into the multiplex format, or demultiplexing by separating or dropping out signals, the filters and couplers used must have center wavelength stability, with a high degree of precision. This in turn means that temperature variations cannot be permitted to affect operation over a substantial range, such as 0-70.degree. C. For a number of reasons, Bragg grating devices are widely employed in wavelength division multiplex systems because a precise wavelength response can be set by selection of the periodicity of the grating written in a device. Because the periodicity of the grating is affected by temperature, however, due to the temperature coefficient of the waveguide or fiber material, the wavelength selectivity can be unacceptably changed by ambient conditions. For example, without any form of temperature compensation, the center wavelength can vary by as much as +0.01 nm/C. Channel placements and bandwidths for individual signals are such that this amount of variation of length, in a normal temperature environment, would shift the wavelength selectivity such that the coupler or filter is no longer matched to the wavelength of the corresponding laser or other source.\nIn a prior application, Ser. No. 09/128,476 entitled \"Methods of Fabricating Grating Assisted Coupler Devices\", filed Aug. 4, 1998 by Anthony S. Kewitsch, et al., and assigned to the assignee of the present invention, this problem is confronted by a disclosed temperature compensation system in which a \"prepackage\" or support structure is disposed within a cylindrical housing, with the prepackage structure using a number of spaced apart hubs mounted on and/or selectively movable relative to Invar rods. The prepackage structure enables attachment of a span of a fiber-based coupler, and the arrangement provides for the use of thermally dissimilar metal elements (stainless steel) and adjustment means for varying the strain on the critical fiber length which includes a Bragg grating. While this arrangement is fully satisfactory from the operational standpoint, and provides for the needed tensioning of the fiber optic components during the initial writing phase and thereafter despite temperature variations, it is more complex than desirable for high production processes. Further, it is desirable to provide even better temperature compensation, to a level of + or -0.001 nm/.degree. C. center wavelength variation over temperature. In addition, it is desirable to provide means for axial twisting of the fiber optic component, for purposes of achieving polarization independence and reducing polarization mode dispersion (PMD)."}
{"text": "In optical communication networks using optical fibers, an increase in the amount of communication data has pushed up the data rate. Optical signals carrying high-speed data over, for example, 10 Gb/s, which are communicated via optical communication networks, are generated by driving optical modulators by electrical signals including the above data. As an example, a data signal that is an electrical signal for driving an optical modulator having a Mach-Zehnder configuration is generated through parallel-to-serial conversion performed on a plurality of parallel data signals to produce a single serial data signal. FIG. 6 illustrates a parallel-to-serial converter, and FIG. 7 illustrates a timing chart of the parallel-to-serial converter. The circuit illustrated in FIG. 6 is illustrated in, for example, “Design of Integrated Circuits for Optical Communications”, Behzad Razavi, International Edition 2003, (Singapore), McGraw-Hill Education, 2003, pp. 333-339.\nThe parallel-to-serial converter in FIG. 6 is a circuit that converts two input data signals IN0 and IN1, which are input in parallel, into a single serial signal to generate an output data signal OUT whose data rate is twice that of the input signals.\nThe input data signal IN0 is latched sequentially by three cascade-connected flip-flops 1-1, 1-2, and 1-3, and is input to a selection circuit (selector) 2 as a selection data signal D0. The input data signal IN1 is latched sequentially by two cascade-connected flip-flops 1-4 and 1-5, and is input to the selection circuit 2 as a selection data signal D1. Each of the flip-flops 1-1 to 1-5 serving as latch circuits operates in accordance with a frequency-divided clock signal divCLK that is obtained using a frequency divider circuit 3 by dividing the frequency of a clock signal CLK whose frequency corresponds to the data rate of the output data signal OUT by ½. The flip-flops 1-2 and 1-5 execute setup and hold operations in accordance with a rising edge of the frequency-divided clock signal divCLK, and the flip-flops 1-1, 1-3, and 1-4 execute setup and hold operations in accordance with a falling edge of the frequency-divided clock signal divCLK.\nTherefore, the input data signals IN0 and IN1 are captured into the flip-flops 1-1 and 1-4, respectively, in accordance with a falling edge of the frequency-divided clock signal divCLK. The input data signals IN0 and IN1 are sequentially sent to the subsequent flip-flops 1-2, 1-3, and 1-5 for every half period of the frequency-divided clock signal divCLK. The selection data signal D0 output through the three flip-flops 1-1, 1-2, and 1-3, and the selection data signal D1 output through the two flip-flops 1-4 and 1-5 are input to the selection circuit 2 as signals whose phases are different by 180° from each other. In this manner, the timing at which the selection data signals D0 and D1 are input to the selection circuit 2 is adjusted using the flip-flops 1-1 to 1-5 that are included in an input circuit.\nThe frequency-divided clock signal divCLK is supplied to the selection circuit 2 as a switching clock signal. The selection circuit 2 performs a switching operation in accordance with the frequency-divided clock signal divCLK to alternately select the selection data signals D0 and D1 to perform serial conversion to produce a conversion data signal D[0+1], and outputs the conversion data signal D[0+1]. The selection circuit 2 selects the selection data signal D0 in accordance with a rising edge of the frequency-divided clock signal divCLK, and selects the selection data signal D1 in accordance with a falling edge of the frequency-divided clock signal divCLK. Therefore, the conversion data signal D[0+1] output from the selection circuit 2 is output as a serial data signal whose frequency corresponds to the frequency of the clock signal CLK whose frequency has not been divided.\nThe conversion data signal D[0+1] is latched in a flip-flop 4 that performs a latch operation in accordance with an output clock signal ffCLK having substantially the same cycle as the clock signal CLK, and the output data signal OUT is output from the flip-flop 4. The flip-flop 4 serving as an output circuit executes setup and hold operations in accordance with a rising edge of the output clock signal ffCLK.\nAs illustrated in the timing chart of FIG. 7, the conversion data signal D[0+1] output from the selection circuit 2 is a data signal in which the selection data signal D0 and the selection data signal D1 are alternately switched with the width of one pulse of the frequency-divided clock signal divCLK. Further, the phases of the clock signals ffCLK and divCLK, which are generated based on the clock signal CLK, are designed so that a rising edge of the output clock signal ffCLK whose frequency is twice that of the frequency-divided clock signal divCLK is located at the center of each pulse of the frequency-divided clock signal divCLK. This enables the flip-flop 4 to accurately latch data at optimum points (the center of pulses) of the conversion data signal D[0+1] and to output a normal output data signal OUT.\nHowever, a phase error with respect to the design, that is, a phase deviation, may occur between the frequency-divided clock signal divCLK and the output clock signal ffCLK because the frequency-divided clock signal divCLK output from the frequency divider circuit 3 is supplied to a large number of elements through a buffer BUF, because there is a difference between wiring capacitances caused by the different wiring lengths of the output clock signal ffCLK and the frequency-divided clock signal divCLK, and other reasons.\nIf a phase deviation occurs between the output clock signal ffCLK supplied to the flip-flop 4 and the frequency-divided clock signal divCLK for outputting the conversion data signal D[0+1], accurate data is not output. This will be described with reference to FIGS. 8A and 8B. FIGS. 8A and 8B are timing charts illustrating the waveform of the conversion data signal D[0+1] and the waveform of the output clock signal ffCLK. FIG. 8A illustrates the waveforms under normal conditions, and FIG. 8B illustrates the waveforms when a phase deviation occurs.\nThe conversion data signal D[0+1] is output for a cycle corresponding to the pulse width of one half period of the frequency-divided clock signal divCLK. Thus, if a rising edge of the output clock signal ffCLK is within a given margin (allowed range) according to the pulse width, normal data is latched in the flip-flop 4 (FIG. 8A). However, if a phase deviation occurs between the frequency-divided clock signal divCLK and the output clock signal ffCLK because of the above reasons and abnormal data is latched in the flip-flop 4 (FIG. 8B) due to this phase deviation, a rising edge of the output clock signal ffCLK is out of the margin.\nAs described above, with an increase in data rate, the rate of the clock signal CLK increases and the pulse width of the frequency-divided clock signal divCLK and the output clock signal ffCLK, which are generated based on the clock signal CLK, is reduced. Then, the margins illustrated in FIGS. 8A and 8B become small. It is therefore desirable that the timing at which the above clock signals are input to the respective circuits be set accurately. A technology for providing a replica circuit to facilitate phase matching is proposed in, for example, “A Single-40 Gb/s Dual-20 Gb/s Serializer IC with SFI5.2 Interface in 65 nm CMOS”, Kouichi Kanda, et al., ISSCC 2009/SESSION 21/10 Gb/s-T0-40 Gb/s TRANSMITTERS AND RECEIVERS/21.2, 2009 IEEE International Solid-State Circuits Conference. FIG. 9 illustrates a parallel-to-serial converter including the above replica circuit.\nIn the parallel-to-serial converter illustrated in FIG. 9, a conversion circuit for performing parallel-to-serial conversion on input data signals IN0 and IN1 to generate an output data signal OUT (that is, flip-flops 1-1 to 1-5 in an input circuit, a selection circuit 2, a frequency divider circuit 3, and a flip-flop 4 in an output circuit) is substantially the same as that in the parallel-to-serial converter illustrated in FIG. 6. A clock signal CLK, a frequency-divided clock signal divCLK, and an output clock signal ffCLK are also substantially the same as those in FIG. 6.\nThe replica circuit includes a selection circuit (replica selection circuit) 10 that selects one of replica input data signals repIN0 and repIN1 in accordance with the frequency-divided clock signal divCLK, and a flip-flop 11 in a replica output circuit that latches a replica conversion data signal repD, which is obtained through serial conversion and output from the selection circuit 10, in accordance with a replica clock signal repCLK.\nThe selection circuit 10 has substantially the same structure as the selection circuit 2, and substantially the same signal as the frequency-divided clock signal divCLK supplied to the selection circuit 2 is supplied as a switching clock signal. Further, the flip-flop 11 also has the same structure as the flip-flop 4. Like the output clock signal ffCLK supplied to the flip-flop 4, the replica clock signal repCLK supplied to the flip-flop 11 is a signal having substantially the same cycle as the clock signal CLK. However, the flip-flop 11 executes setup and hold operations in accordance with a falling edge of the replica clock signal repCLK. In this manner, the selection circuit 10 and the flip-flop 11 are replicas of the selection circuit 2 and the flip-flop 4, respectively.\nIn the replica circuit, a replica output data signal repOUT output from the flip-flop 11 is latched in accordance with a latch clock signal latCLK output from the frequency divider circuit 3. The phase of the frequency-divided clock signal generated from the frequency divider circuit 3 is adjusted using this latched value.\nThe replica output data signal repOUT output from the flip-flop 11 is captured into a flip-flop 12 provided as a latch circuit. The flip-flop 12 performs a latch operation in accordance with the latch clock signal latCLK output from the frequency divider circuit 3 whose frequency has been divided by half the frequency of the clock signal CLK. Then, the flip-flop 12 executes setup and hold operations in accordance with a rising edge of the latch clock signal latCLK. The data latched in the flip-flop 12 is supplied to a counter 13 (see “A Single-40 Gb/s Dual-20 Gb/s Serializer IC with SFI5.2 Interface in 65 nm CMOS”, Kouichi Kanda, et al., ISSCC 2009/SESSION 21/10 Gb/s-T0-40 Gb/s TRANSMITTERS AND RECEIVERS/21.2, 2009 IEEE International Solid-State Circuits Conference; FIG. 21.2.3, U/D counter). The counter 13 counts up or counts down in accordance with the output value of the flip-flop 12 in synchronization with the latch clock signal latCLK produced by the frequency divider circuit 3. For example, in accordance with a rising edge of the latch clock signal latCLK, the count value is increased by 1 if the current output value of the flip-flop 12 is “1” (high level), and the count value is decreased by 1 if the output value is “0” (low level). The flip-flop 12 serving as a latch circuit and the counter 13 form a phase deviation detection circuit that generates a phase adjustment instruction value signal based on the replica output data signal repOUT.\nThe count value output from the counter 13 as a phase adjustment instruction value signal is decoded at a desired time using a decoder 14, and is converted into an appropriate code (for example, a thermometer code). Based on the code, an arbitrary phase generation circuit 15 adjusts the phase of the frequency-divided clock signal output from the frequency divider circuit 3 (see “A Single-40 Gb/s Dual-20 Gb/s Serializer IC with SFI5.2 Interface in 65 nm CMOS”, Kouichi Kanda, et al., ISSCC 2009/SESSION 21/10 Gb/s-T0-40 Gb/s TRANSMITTERS AND RECEIVERS/21.2, 2009 IEEE International Solid-State Circuits Conference; FIG. 21.2.3, PI: Phase Interpolator), and generates a frequency-divided clock signal divCLK. The arbitrary phase generation circuit 15 serving as a phase interpolation circuit outputs a frequency-divided clock signal divCLK whose phase has been adjusted, by, for example, weighting multiphase frequency-divided clock signals having substantially the same cycle as the frequency-divided clock signal output from the frequency divider circuit 3 and having phases different by 90° from each other in accordance with the code supplied from the decoder 14 and superposing the weighted clock signals.\nFIG. 10 illustrates a timing chart of the parallel-to-serial converter in FIG. 9.\nThe replica input data signals repIN0 and repIN1 are input as data signals one of which is fixed to “0” and the other of which is fixed to “1”. Therefore, the replica conversion data signal repD output from the selection circuit 10 transitions between “0” and “1” every half period of the frequency-divided clock signal divCLK. The flip-flop 11 performs a latch operation in accordance with a falling edge of the replica clock signal repCLK whose frequency is twice that of the frequency-divided clock signal divCLK. Thus, the flip-flop 11 executes a latch for the transition period of the replica conversion data signal repD, and generates a replica output data signal repOUT. Therefore, in the replica output data signal repOUT that is latched in and output from the flip-flop 11, “0” and “1” appear randomly. Specifically, a fine change of the transition timing of the replica conversion data signal repD with respect to the latch timing of the flip-flop 11 and a fine change of the signal level of the replica conversion data signal repD at the latch timing cause the replica output data signal repOUT not to appear in a predictable pattern. That is, “0” and “1” appear with substantially the same probability.\nThe flip-flop 12 that latches the replica output data signal repOUT operates in accordance with the latch clock signal latCLK output from the frequency divider circuit 3. The latch clock signal latCLK is a signal that has just been output from the frequency divider circuit 3, that is, a signal that has not yet passed through the arbitrary phase generation circuit 15 or the buffer BUF, and the phase does not substantially change with respect to the clock signal CLK (that is, the replica clock signal repCLK). The flip-flop 12 that operates in accordance with the latch clock signal latCLK is allowed to execute a latch at any timing out of the transition period of the replica output data signal repOUT. Since, as described above, the replica output data signal repOUT latched in the flip-flop 12 is a data signal in which “0” and “1” appear randomly, the flip-flop 12 may latch “0” or “1” with substantially the same probability.\nAs a result of “0” or “1” being latched in the flip-flop 12 with substantially the same probability, the count value of the counter 13 that counts the value “0” or “1” is maintained at a substantially constant value and is made stable although the count value may go up and down slightly. In accordance with the stable count value, a decoding value output from the decoder 14 is also made stable, and the frequency-divided clock signal divCLK output from the arbitrary phase generation circuit 15 in accordance with the decoding value is made stable.\nFIGS. 11A and 11B illustrate timing charts when a phase deviation occurs between the frequency-divided clock signal divCLK and the replica clock signal repCLK (that is, the output clock signal ffCLK) due to the reasons described above. FIG. 11A illustrates a result obtained when the phase of the replica clock signal repCLK is delayed with respect to the frequency-divided clock signal divCLK, and FIG. 11B illustrates a result obtained when the phase of the replica clock signal repCLK is advanced with respect to the frequency-divided clock signal divCLK. The output clock signal ffCLK and the replica clock signal repCLK may be regarded as the same signal based on the clock signal CLK. Therefore, the replica clock signal repCLK is illustrated in FIGS. 11A and 11B, by way of example.\nAs described above, there is substantially no phase deviation between the replica clock signal repCLK and the latch clock signal latCLK. Thus, in FIG. 11A, the phase of the frequency-divided clock signal divCLK may be regarded as being deviated ahead of the replica clock signal repCLK and the latch clock signal latCLK. The replica conversion data signal repD transitions in response to the frequency-divided clock signal divCLK while the flip-flop 11 performs a latch operation in response to the replica clock signal repCLK. Thus, if the phases of the clock signals divCLK and repCLK are deviated from each other, a deviation may occur in the latch timing of the flip-flop 11 with respect to the transition period of the replica conversion data signal repD. If the latch timing is not within the transition period, the values latched in the flip-flop 11 do not appear randomly but repetitions of “0” and “1” according to the replica conversion data signal repD may occur. Therefore, the replica output data signal repOUT may be a data signal in which “0” and “1” repeatedly appear every period of the replica clock signal repCLK.\nThe flip-flop 12 that performs a latch operation in response to the latch clock signal latCLK in which substantially no phase deviation occurs from the replica clock signal repCLK captures the replica output data signal repOUT on every other pulse because the latch clock signal latCLK has a cycle that is twice as that of the replica clock signal repCLK. In FIG. 11A, only the values “0” of the replica output data signal repOUT are latched. Consequently, the output of the flip-flop 12, which is indicated by “FF12” in FIG. 11A, is maintained at “0”. The counter 13 decreases the count value in accordance with this output. Thus, the arbitrary phase generation circuit 15 that operates in accordance with the count value and the decoding value delays the phase of the frequency-divided clock signal divCLK. Finally, the timing of each clock signal is adjusted to an optimum timing in FIG. 10.\nIn FIG. 11B, the phase of the frequency-divided clock signal divCLK may be regarded as being deviated and delayed with respect to the replica clock signal repCLK and the latch clock signal latCLK. In this case, the flip-flop 12 in FIG. 11B, which captures the replica output data signal repOUT on every other pulse, latches only the values “1” of the replica output data signal repOUT. Consequently, the output of the flip-flop 12, which is indicated by “FF12” in FIG. 11B, is maintained at “1”. The counter 13 increases the count value in accordance with this output. Thus, the arbitrary phase generation circuit 15 that operates in accordance with the count value and the decoding value advances the phase of the frequency-divided clock signal divCLK. Finally, the timing of each clock signal is adjusted to an optimum timing in FIG. 10.\nIn some cases, a parallel-to-serial converter including the replica circuit described above may be provided with a plurality of conversion circuits each of which performs parallel-to-serial conversion on a data signal, that is, a plurality of conversion circuits each including the flip-flops 1-1 to 1-5 in the input circuit, the selection circuit 2, and the flip-flop 4 in the output circuit. In such cases, the plurality of conversion circuits are driven simultaneously through clock phase adjustment using a single replica circuit. FIG. 12 illustrates this parallel-to-serial converter.\nThe clock signals used in the parallel-to-serial converter in FIG. 12 are similar to the corresponding clock signals in FIG. 9. Further, a replica circuit including a selection circuit 10, flip-flops 11 and 12, a counter 13, a decoder 14, and an arbitrary phase generation circuit 15, and a frequency divider circuit 3 are also substantially the same as those in FIG. 9.\nIn FIG. 12, two conversion circuits each of which converts a parallel data signal into a serial data signal are provided. A first conversion circuit includes flip-flops 1-1a to 1-5a in an input circuit that latch input data signals IN0a and IN1a in accordance with a frequency-divided clock signal divCLK produced by the arbitrary phase generation circuit 15, a selection circuit 2a that selects one of selection data signals D0a and D1a in accordance with the frequency-divided clock signal divCLK, and a flip-flop 4a in an output circuit that latches a conversion data signal D[0+1]a in accordance with an output clock signal ffCLK and that outputs an output data signal OUTa. A second conversion circuit has substantially the same structure, and includes flip-flops 1-1b to 1-5b in an input circuit that latch input data signals IN0b and IN1b in accordance with the frequency-divided clock signal divCLK, a selection circuit 2b that selects one of selection data signals D0b and D1b in accordance with the frequency-divided clock signal divCLK, and a flip-flop 4b in an output circuit that latches a conversion data signal D[0+1]b in accordance with the output clock signal ffCLK and that outputs an output data signal OUTb.\nThat is, the parallel-to-serial converter in FIG. 12 is provided with a plurality of conversion circuits having substantially the same structure, and each conversion circuit is driven by commonly using the frequency-divided clock signal divCLK produced by the frequency divider circuit 3 and the arbitrary phase generation circuit 15 and the output clock signal ffCLK produced from the clock signal CLK. The common frequency-divided clock signal divCLK is phase-adjusted using a single replica circuit.\nThe replica clock signal repCLK supplied to the flip-flop 11 in the replica circuit and the output clock signal ffCLK supplied to the flip-flops 4a and 4b in the respective conversion circuits are signals that are divided from the same signal, namely, the clock signal CLK. Preferably, the replica clock signal repCLK and the output clock signal ffCLK are delivered to the flip-flops 4a, 4b, and 11 using isometric lines having a substantially equal wiring capacitance to prevent a phase deviation from occurring therebetween. However, as in FIG. 12, if a plurality of conversion circuits are provided, due to the layout design on an integrated circuit (IC) chip, it may be difficult to make the wiring length for the replica clock signal repCLK directed to the replica circuit equal to the wiring length for the output clock signal ffCLK directed to the conversion circuits. In addition, even when it is possible to make the lengths of the wirings substantially equal to each other, with the continued miniaturization of circuits, the influence of the difference between the input capacitances of flip-flops on the phase of clocks has increased.\nA phase deviation between the replica clock signal repCLK and the output clock signal ffCLK will be described with reference to a timing chart of FIG. 13.\nA replica circuit is a circuit that adjusts the phase of the frequency-divided clock signal divCLK with respect to the replica clock signal repCLK, but does not adjust the phase of the frequency-divided clock signal divCLK with respect to the output clock signal ffCLK. That is, the replica circuit is configured based on the assumption that there is no phase deviation between the replica clock signal repCLK and the output clock signal ffCLK that are generated from substantially the same signal, namely, the clock signal CLK. A phase deviation between the replica clock signal repCLK and the output clock signal ffCLK due to the reasons described above may cause the influence of the output clock signal ffCLK on the latch timing of the flip-flops 4a and 4b. \nIn the timing chart of FIG. 13, due to the reasons described above, the phase of the output clock signal ffCLK is delayed with respect to the replica clock signal repCLK. The phase of the frequency-divided clock signal divCLK is appropriately adjusted with respect to the replica clock signal repCLK using the replica circuit, and the conversion data signals D[0+1]a and D[0+1]b are output at appropriate timings using the selection circuits 2a and 2b in accordance with the frequency-divided clock signal divCLK. Meanwhile, the output clock signal ffCLK is delayed with respect to the replica clock signal repCLK because of the occurrence of a phase deviation. Thus, a rising edge of the output clock signal ffCLK is out of an appropriate timing for latching the conversion data signals D[0+1]a and D[0+1]b, that is, the margins illustrated in FIGS. 8A and 8B. This may cause, as described with reference to FIG. 8B, the flip-flops 4a and 4b to latch incorrect data, and the output data signals OUTa and OUTb, which are abnormal, are output.\nA parallel data output device illustrated in FIG. 14 as well as a parallel-to-serial converter may also experience difficulties involved in a phase deviation of clock signals. The parallel data output device in FIG. 14 is provided to, for example, drive an optical modulator OM used in multilevel phase modulation methods. The optical modulator OM is an optical component that performs multilevel phase modulation of continuous light emitted from a light source LS using a Mach-Zehnder optical modulation circuit and that outputs a modulated optical signal. Parallel output data signals OUT0 and OUT1 that are produced by a parallel data output device are used as drive signals to be applied to a traveling wave electrode of the Mach-Zehnder optical modulation circuit. It is therefore desirable that the output data signals OUT0 and OUT1 produced by the parallel data output device be synchronous signals whose timings are made to match with high accuracy.\nThe parallel data output device is configured using flip-flops 20 and 21 serving as latch circuits that latch the input data signals IN0 and IN1 input thereto in parallel, respectively. The flip-flops 20 and 21 perform a latch operation in accordance with a rising edge of output clock signals ffCLK0 and ffCLK1 that are distributed from a single clock signal CLK, thereby outputting the output data signals OUT0 and OUT1 at a synchronous timing.\nSince the output clock signals ffCLK0 and ffCLK1 are generated from substantially the same clock signal, namely, the clock signal CLK, a phase deviation is unlikely to occur between these signals. However, as described above, the influence of the difference between the input capacitances of the flip-flops 20 and 21 has increased in accordance with the continued miniaturization of circuits, resulting in a phase deviation being likely to occur between the output clock signals ffCLK0 and ffCLK1. A phase deviation between the output clock signals ffCLK0 and ffCLK1 may influence the deviation between the output timings of the output data signals OUT0 and OUT1. The deviation between the output timings of the output data signals OUT0 and OUT1 is illustrated in a timing chart of FIG. 15."}
{"text": "This invention relates to 2-(oxa and thia heterocycle) 5-amino-3-oxo-4-(substituted-phenyl)-2,3-dihydrofurans and derivatives thereof and to the use of such compounds as herbicides and plant growth regulators.\nChemiker-Zeitung 104 (1980) No. 10, Pages 302-303, is an academic paper disclosing the ring closure of 1-(dimethylamino)-2,4-diphenyl-1-buten-3,4-dione to yield 5-dimethylamino-2,4-diphenyl-2,3-dihydrofuran. British Pat. No. 1,521,092, discloses certain 3-phenyl-5-substituted-4(1H)-pyrid-ones or -thiones as herbicides. Japanese Patent Application No. 13,710/69 (Chemical Abstracts 71:61195e) discloses 5-amino-3-oxo-4-(phenyl or 4-chlorophenyl)-2,3-dihydrofurans. Japanese Patent No. 19090 (Chemical Abstracts 69P10352e) discloses certain 2,3-dihydrothiophenes as pharmaceuticals. Helvetica Chemica Acta, Volume 66, Pages 362-378 (1983) discloses 5-N-cyclopropyl-4-phenyl-2-methoxycarbonylmethylene-3-furanone as part of an academic chemical synthesis discussion.\nIn my copending application U.S. Ser. No. 505,169, filed June 17, 1983, I disclosed certain 2-aryl 5-amino-3-oxo-4-(substituted phenyl)-2,3-dihydrofuran derivatives and their use as herbicides."}
{"text": "1. Field of the Invention\nThe present invention relates to environment recognition systems, and more particularly, to an environment recognition system that recognizes a surrounding environment by detecting objects in a taken image.\n2. Description of the Related Art\nIn general, in order to measure the distance to an object existing in a surrounding environment with a stereo camera, a pair of images are taken by a pair of right and left cameras that are mounted at the same height, and one of the taken images used for reference (hereinafter referred to as a reference image T0) is compared with the other image (hereinafter referred to as a comparative image Tc). By comparison, a difference between corresponding positions of the same object in the images, that is, a parallax is calculated, and the distance to the object is calculated from the parallax. The positions in the reference image and the comparative image where an image of the same object is included are typically located by stereo matching (for example, see Japanese Unexamined Patent Application Publication Nos. 10-283461 and 10-283477).\nIn stereo matching, as shown in FIG. 17, a reference image T0 is divided into small regions (hereinafter referred to as reference pixel blocks PB0) each defined by a predetermined number of pixels, such as 3 by 3 pixels or 4 by 4 pixels. An epipolar line EPL is set in the vertical position in the comparative image Tc corresponding to each reference pixel block PB0, and a brightness pattern of the reference pixel block PB0 is compared with a brightness pattern of a comparative pixel block PBc that exists on the epipolar line EPL and that has the same shape as that of the reference pixel block PB0.\nIn this case, for example, a SAD (Sum of Absolute Difference) value is calculated as a difference in the brightness pattern according to the following Expression (1):\n                    SAD        =                              ∑                          s              ,              t                                ⁢                                                                p                ⁢                                                                  ⁢                1                ⁢                st                            -                              p                ⁢                                                                  ⁢                2                ⁢                st                                                                                    (        1        )            where p1st represents the brightness of the pixel in the reference pixel block PB0 in the reference image T0, and p2st represents the brightness of the pixel in the comparative pixel block PBc in the comparative image Tc. Of SAD values that are less than or equal to a preset threshold value, a comparative pixel block PBc that provides the smallest SAD value is specified as a comparative pixel block in the comparative image Tc that includes an image of the same object as that included in the reference pixel block PB0.\nA parallax dp between the comparative pixel block PBc specified in the comparative image Tc and the original reference pixel block PB0 in the reference image T0 is calculated, and a distance Z to the object at the reference pixel block PB0 is calculated on the basis of the parallax dp according to the principle of triangulation. On the basis of the calculated distance Z, the object is detected from the surrounding environment.\nIt is confirmed that this object detection method that calculates the parallax dp by stereo matching of the reference image T0 and the comparative image Tc and calculates the distance Z to the object functions without any trouble in a normal image taking environment and can effectively detect the object from the surrounding environment, as disclosed in the above-described publications.\nHowever, for example, when the stereo camera is placed in a backlit environment, a reference image T0 shown in FIG. 18A is bright because backlight enters the image, while a comparative image Tc shown in FIG. 18B taken in the same scene is totally darker than the reference image T0 because backlight is blocked by a building or the like and much backlight does not enter the image.\nWhen the brightness balance between a pair of cameras is thus disturbed, the difference between the brightness p1st of the pixel in the reference pixel block PB0 in the reference image T0 and the brightness p2st of the pixel in the comparative pixel block PBc in the comparative image Tc in Expression (1) described above generally increases. Therefore, the calculated SAD value increases above the above-described threshold value. In this case, the parallax dp is not effectively calculated, and the number of reference pixel blocks PB0 to be rejected increases.\nFor this reason, in data image (hereinafter referred to as a distance image Tz) formed by assigning calculated parallaxes dp to the pixel blocks PB0 in the reference image T0, little data on effective parallaxes dp is provided, as shown in FIG. 19. In this case, it is sometimes difficult to detect objects, and reliability of the object detection result decreases. In the worst case, little data on effective parallaxes dp is included in the obtained distance image Tz, and it is completely impossible to detect objects.\nIn this case, for example, a reference edge image TE0 shown in FIG. 20A is formed by calculating differences in the brightness p1ij between pixels belonging to the reference image T0 shown in FIG. 18A and pixels adjacent on the right or left side. Similarly, a comparative edge image TEc shown in FIG. 20B is formed from the comparative image Tc shown in FIG. 18B. The edge image TE0 and the comparative edge image TEc can be subjected to stereo matching.\nBy subjecting the reference edge image TE0 and the comparative edge image TEc thus formed to stereo matching, a distance image TEz in which a relatively large amount of data on effective parallaxes dp are included (hereinafter a distance image based on the edge images is referred to as an edge distance image) is obtained, as shown in FIG. 21. In this case, objects can sometimes be effectively detected even when it is difficult to effectively detect the objects by directly conducting stereo matching on the original reference and comparative images T0 and Tc. FIGS. 20A, 20B, and 21 illustrate parts of the reference edge image TE0, the comparative edge image TEc, and the edge distance image TEz.\nHowever, edge processing has a problem in that much information is lost when obtaining the differences in brightness between the adjacent pixels. That is, when the difference in brightness is 30 in 256 brightness levels, it is unclear whether the difference of 30 indicates a difference between 50 and 80 or between 200 and 230. Moreover, this amplifies noise components in the reference image T0 and the comparative image Tc.\nFurther, since the difference in brightness only in a width corresponding to one or several pixels is found, information about low-frequency components in the frequency components in the reference image T0 and the comparative image Tc is lost. Therefore, mismatching easily occurs. Further, it is difficult to obtain effective information, for example, about a wall that is not characteristic in structure and pattern and an asphalt road surface.\nIn this way, although edge processing is effective, as described above, it should be avoided to always detect objects only on the basis of an edge distance image TEz that is formed by stereo matching of a reference edge image TE0 and a comparative edge image TEc obtained by subjecting a reference image T0 and a comparative image Tc to edge processing.\nWhen objects are detected only on the basis of the distance image Tz obtained from the reference image T0 and the comparative image Tc, objects can be effectively and accurately detected from the surrounding environment in a normal image taking condition, as described above. However, it is undeniable that object detection is difficult in the above-described special condition."}
{"text": "1. Field of the Invention\nThis invention relates generally to actuators for ergonomic systems and, more particularly, to actuator systems for seat adjustments.\n2. Related Art\nErgonomic supports for seats, such as lumbar and bolster systems, are typically moved by means of actuators that can be operated by hand or driven by a motor. Four-way power lumbar devices, with an arching mode and a translation mode, have traditionally required a motor for each mode of operation. In comparison, the invention set forth by U.S. Pat. No. 6,050,641 is a four-way power lumbar system that requires only a single-motor and reduces the duplication of gearbox components. Prior to this invention, U.S. Pat. Nos. 5,197,780 and 5,217,278 had only described four-way lumbar devices that were manually operated by a single control knob.\nWhile these devices are an improvement over the conventional four-way lumbar devices that required multiple motors or multiple control knobs, the highly competitive markets for furniture and automotive seats place a premium on continued optimization of devices that provide comfort and convenience for seat occupants. In particular, there is a need for improved actuation systems that are less prone to failures and more efficiently transfer power to the actuators. For example, with regard to the manually-operated four-way lumbar systems, the gearing systems are inefficient because, in addition to the gears, they have a transmission that requires at least one non-gear alignment device to maintain the proper engagement between the driver gear and the driven gears. Therefore, such devices have an alignment device that is necessary for the convenient operation of the lumbar system and increase the potential for a failure in the system.\nAdditionally, there is a need for actuation systems that are more modular, increasing the commonality of parts between two-way and four-way power lumbar devices, manually-operated and motor-driven actuation systems, and lumbar systems and bolster systems. For example, the prior art actuation units and driven gears are designed to fit within a single housing along with the driver gears and the transmission system, thereby limiting the range of motion that is capable for the actuation units themselves. Different types of lumbar devices and bolster devices are typically designed to provide different levels of support and often require different levels of actuation, thereby affecting the size of the actuators. The prior art devices do not easily allow for changing the actuators according to various sizing requirements, and the confined housing could prevent the same actuation system from being used for different lumbar devices or for a lumbar device and a bolster device. Therefore, entirely different actuation systems would need to be designed, and separately manufactured, depending on the actuators' usage, range of motion, and sizing."}
{"text": "As semiconductor fabrication technologies are continually progressing to smaller feature sizes such as 65 nanometers, 45 nanometers, and below, immersion lithography methods are being adopted. However, during an exposure process using an immersion lithography system, contaminations such as bubbles, particles and water residues can be introduced into the immersion lithography system and further contaminate semiconductor wafers to be processed therein. Such contamination can cause defects and yield degradations."}
{"text": "A technology has been conventionally known that reduces the amount of energy consumption of a service side by altering a service timetable so as not to exceed a certain amount of electric power that is under contract with an electric power company. In such a technology, for the purpose of energy saving of a whole of a railroad system including station facilities, the energy saving of a station side may also be performed, in addition to the energy saving on the service side. In this case, it is common that the energy saving on the service side and the energy saving on the station side are performed independently from each other.\nIn such a technology described above, for an example, it is desirable that the energy saving on the service side and the energy saving on the station side be able to be performed in a unified manner."}
{"text": "A technique has hitherto been known in which, in order to wake a subject from a short sleep completely, a sleep rhythm is estimated from the amount of sweating of the subject, and after the subject has once reached a REM (Rapid Eye Movement) sleep, a stimulator is activated to wake the subject (for example, see Patent Literature 1)."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for manufacturing a polyurethane (P.U.) product with high wetting ability, gas permeability and high water repellent ability by means of a dry transfer coating process.\n2. Description of the Related Art\nA conventional dry type direct coating method in accordance with the prior art shown in FIGS. 1 and 2 is used for manufacturing a synthetic cloth and comprises directly coating a P.U. surface layer 20 on a substrate cloth 10 which is a knitting cloth, a non-woven cloth, intersecting knitting cloth or the like, thereby making the synthetic cloth which includes the P.U. surface layer 20 and the substrate cloth 10 as shown in FIG. 2. In such a manner, a film layer with small apertures is formed on the substrate cloth 10 for increasing the gas permeability of the synthetic cloth, wherein the wetting ability of the synthetic cloth is about 1500 to 2000 g/SQMD. However, the P.U. surface layer 20 cannot be formed with proper patterns so as to fit the different requirements of the consumers, thereby limiting the versatility of the synthetic cloth.\nA conventional wet type transfer coating method in accordance with the prior art shown in FIGS. 3 and 4 is used for manufacturing a synthetic cloth and comprises the following steps. First, a P.U. surface layer 20 is coated on a releasing sheet 1, and is then baked and dried. A bonding agent 30 is then bonded on the P.U. surface layer 20 for bonding a substrate cloth 10 to the P.U. surface layer 20. The substrate cloth 10 is a knitting cloth, a non-woven cloth, intersecting knitting cloth or the like. The bonding agent 30 is then baked and dried, and the releasing sheet 1 is then removed from the P.U. surface layer 20, thereby forming the final synthetic cloth which includes the P.U. surface layer 20, the bonding agent 30, and the substrate cloth 10 as shown in FIG. 4. A plurality of small apertures are formed in the P.U. surface layer 20 for increasing the gas permeability of the synthetic cloth, wherein the wetting ability of the synthetic cloth is decreased to about 650 g/SQMD. In such a manner, the patterns on the releasing sheet 1 is printed on the P.U. surface layer 20 by a transfer type. However, the wetting ability and gas permeability of the synthetic cloth are decreased.\nIn accordance with one aspect of the present invention, there is provided a method for manufacturing a polyurethane (P.U.) product with high wetting ability, gas permeability and high water repellent ability by means of a dry transfer coating process. The method comprises a front stage working, and a rear stage working.\nThe front stage working includes the steps of: providing a releasing sheet which has a plurality of pattern cavities defined therein; providing a first P.U. liquid which is scraping resistant and strong solvent resistant; coating the first P.U. liquid on the releasing sheet to fill the pattern cavities of the releasing sheet with the first P.U. liquid; baking and drying the first P.U. liquid to form a first P.U. surface layer on the releasing sheet; providing a second weak solvent P.U. liquid; coating the second weak solvent P.U. liquid on the first P.U. surface layer; baking and drying the second weak solvent P.U. liquid to form a second P.U. surface layer on the first P.U. surface layer; providing a bonding agent on the second P.U. surface layer; providing a substrate cloth on the bonding agent; and baking and drying the bonding agent so that the substrate cloth is securely bonded on the second P.U. surface layer through the bonding agent, thereby forming a first product which includes the first P.U. surface layer, the second P.U. surface layer, the bonding agent, and the substrate cloth.\nThe rear stage working includes the steps of: providing a first mixing liquid to be rolled and printed on the first product for performing a surface treatment of a first printing plate so as to form a plurality of intensive small air vents in the second P.U. surface layer for increasing permeability of the second P.U. surface layer; providing a second mixing liquid containing a fluoride to be rolled and printed on the first product for performing a surface treatment of a second printing plate so that the fluoride infiltrates through the first P.U. surface layer and the second P.U. surface layer for increasing a water repellent function of the first P.U. surface layer; and providing a third mixing liquid containing a P.U. solvent to be rolled and printed on the first product for performing a surface treatment of a third printing plate so as to increase comfort and quality of the first product, thereby making a final synthetic cloth product.\nA primary, objective of the present invention is to provide a method which uses a transfer coating process to achieve a front stage working to form various patterns. The front stage working includes two coating steps, wherein the first coating step includes coating a scraping resistant and strong solvent resistant first P.U. liquid on the releasing sheet to form a first P.U. surface layer, and the second coating step includes coating a second weak solvent P.U. liquid on the first P.U. surface layer to form a second P.U. surface layer. The method also includes a rear stage working which uses a strong solvent made of D.M.F. to form a plurality of intensive small air vents in the second P.U. surface layer, wherein the transfer patterns printed by the first P.U. surface layer will not be broken.\nAnother objective of the present invention is to enhance the gas permeability and wetting ability of the synthetic cloth product by means of the intensive small air vents formed in the second P.U. surface layer.\nA further objective of the present invention is to enhance the water repellent and waterproof effect of the synthetic cloth product by means of the fluoride.\nA further objective of the present invention is to enhance the quality and comfort of the synthetic cloth product by means of the anti-blocking agent, the surfactant of wax type, and the surfactant of mat and bright."}
{"text": "This invention relates generally to home tools and appliances and more particularly relates to a rain gutter cleaning tool which may be manually held and manipulated by a user for cleaning a rain gutter disposed at a relatively remote location, i.e., a rain gutter situated on the eave of an architectural structure such as a dwelling place, a garage, or a utility building and susceptible to malfunction due to the accumulation of fallen leaves and other yard debris."}
{"text": "The present invention is directed to an insulating foam board system and methods of producing the same. In particular, the present invention is directed to a foam board system and methods thereto that combine at least three layers to provide greater durability and structural integrity to the foam board system.\nInsulating material is used in the construction of buildings. Insulating materials are especially called for in regions of the country that experience extreme weather conditions. Popular modem-day insulating materials include foam boards which are often manufactured from a polystyrene polymer having a laminated coating. The foamed polystyrene board has insulating properties associated therewith. The laminated coating functions primarily to protect the foamed polystyrene polymer and provide the foam board with enhanced protection, durability, strength and resiliency.\nSuch foam board insulation can be found in a number of different forms. For example, foam board insulation is sold in individual flat sheets that may be applied to the exterior of a building in combination with an exterior face material such as brick or siding. Foam board insulation in this application is used as one layer in the overall exterior wall construction of a building.\nFoam board insulation may also be in the form of a folded or hinged board (also referred to as fanfold board) that when unfolded at the hinges can be applied to the exterior of a building in the same manner as flat foam board insulation. The fanfold or flat board insulation can be constructed of a foamed polystyrene with a laminated coating that provides protection, durability, strength and resiliency. The fanfold board is often preferred over the flat foam board insulation because of easier handling and installation.\nNotwithstanding the laminated coating in the above-described types of foam board insulation, the foam board insulation is still susceptible to deterioration at the building site prior to installation. This problem is further evident when, for example, delays occur in the installation and/or the installation occurs in extreme weather conditions. Additionally, the foam board insulation is often mishandled or improperly installed such that the laminated coating does not provide sufficient protection and strength, resulting in damaged foam board insulation. One common example of damage is when a piece breaks off from a remainder of the foam insulation board.\nAccordingly, a need exists for both types of foam board insulationxe2x80x94flat foam board insulation and folded/hinged foam board insulationxe2x80x94which have increased weather resistance, protection, durability and strength.\nIn one embodiment, an insulated foam board system of the present invention comprises a first layer, a second layer and a central layer. The first layer is made from a material selected from the group consisting of alkenyl aromatic polymers, polypropylenes, polyethylene terephthalates, polyethylenes and combinations thereof. The second layer is made from a material selected from the group consisting of alkenyl aromatic polymers, polypropylenes, polyethylene terephthalates, polyethylenes and combinations thereof. The central layer is located between the first layer and the second layer. The central layer is bonded to the first layer and the second layer. The central layer is made from a material selected from the group consisting of crystal polystyrene, impact polystyrene, polyethylene terephthalate and combinations thereof.\nIn another embodiment, the insulated foam board system comprises a first layer, a second layer, a central layer and at least one laminated surface coating."}
{"text": "Membrane separations are a key enabling technology for energy conversion devices. Ionic transport membranes must have both proton and electronic conductivity to function as hydrogen separation membranes without an external power supply. A technical obstacle to material modification by compositional changes is that the hydrogen flux through a dense membrane is a function of both the proton ionic conductivity and the electronic conductivity.\nIn addition, the materials electronic conductivity or material crystal structure stability should not be greatly affected by the presence of contaminant gases such as CO2, CO, CH4 and H2O which are commonly associated with from steam reforming/water gas shift reactions. Perovskite materials of the general formula SrCeO3 and BaCeO3 form the basis of most ceramic compositions with proton conductivities in the range of 2×10−2 S/cm at 600° C., showing good stability under the extremely low oxygen partial pressure where many Perovskites decompose to their primary oxides, “A”-site doping of the ABO3 Perovskite structure and stoichiometry modifications have been explored to increase the stability in the presence of contaminated gases while maintaining acceptable proton conductivity levels.\nThere remains room for variation and improvements in the art of conductive membrane and hydrogen flux through a membrane."}
{"text": "Websites often provide a social sharing feature in the form of a share widget for content hosted by the website. These share widgets enable users to quickly share content with members of their social network using multiple social network services (e.g., Facebook®, Twitter®, Pinterest®, or the like). Although such functionality may provide users with a link between the website's services and the social network services, the websites themselves are often unable to ascertain which of their users correspond to the social network profiles used to publish content from the websites."}
{"text": "The present invention relates to a control and monitoring telecommunication system for connecting and communicating between a site such as power plant and shopping mall where apparatuses to be controlled and monitored are located and a control room and monitor room where operators and observers are stationed, and also to a method of setting a modulation method for the communication.\nA need for remote monitoring or remote maintenance of existing plant has been increasing in many plants including power plant. In a shopping mall or large-scale shop, a need for installing or adding monitor cameras is also increasing in view of higher security.\nIt therefore has become necessary to lay a dedicated cable newly between the site and monitor room or control room so as to transmit the data obtained by the sensors and actuators that have not been connected in a network or monitor cameras, sensors and actuators that have been newly installed from these apparatuses to the monitor room or control room. On the other hand, since there are large apparatuses or inverters installed at the site, eliminating the effect of the electromagnetic noise to be caused by driving these apparatuses needs to be considered. The construction cost needed to lay a dedicated cable in an existing installation so as to avoid the effect or eliminate the noise in the telecommunication turns to be higher than the construction cost of a telecommunication line for measuring and control points spent at the time of initial plant construction.\nGenerally, in a power plant, various sites including turbine room and boiler room are located far from each other by hundreds to thousand meters within an area, and even in a shopping mall, monitor cameras need to be installed as distributed within a relatively wide outdoor space. If a one-to-one communication modem is employed for these applications, a star type network, where individual telecommunication line is laid between each site (monitor camera in the shopping mole application) and the monitor room, is needed but the installation cost and maintenance cost of the total network increases.\nIn order to resolve this problem, there has been proposed a method for utilizing an unused line or frequency band out of the existing telecommunication lines installed between the site and monitor room so as to reduce cost needed to lay a new dedicated line or take an anti-noise measures (see Japanese Laid-open Patent Publication No. 2001-25000).\nAlthough Japanese Laid-open Patent Publication No. 2001-25000 (paragraphs 0026-0035, FIG. 1 and FIG. 2) also discloses a telecommunication mode where multiple modems are connected with a single telecommunication line, it lists nothing else but generally known coordination methods (methods where one modem serving as master allocates the transmission right; comprising four methods: token passing, collision sensing, time division multiplex, and frequency division multiplex). The patent does not contain any concrete description about detailed physical layer of communication but simply refers to a reference literature on xDSL (x digital subscriber line). Detailed description about the physical layer of the x DSL is given in “Draft American National Standard for Telecommunications-Network and Customer Installation Interfaces-Asymmetric Digital Subscriber Line (ADSL) Metallic Interface” (U.S.A.), American National Standards Institute, Inc. ANSI (R) T1. 413-1998, pp. 95-135, for example. The document, which is written on the physical layer of ADSL (asymmetric digital subscriber line), a typical xDSL, describes that, for the modems communicating with each other, training to cope with noise or attenuation and adjusting the communication setting based on the training (determining a modulation method for every carrier) are needed, explaining how to realize it. In this explanation, carrier means the frequency of carrier wave signal."}
{"text": "In a conventional electrolytic machine, liquid electrolyte such as sodium nitrate solution or sodium chloride solution is continuously passed between an electrode and a work at a high speed during machining, so that residual products such as particles of eroded metal from the work, hydrogen gas, and others are discharged from the gap between the electrode and the workpiece and a direct current voltage in the form of pulses is applied between the workpiece and the electrode for finishing the surface of the workpiece.\nThe applicant of the present invention have proposed electrolytic finishing systems in which a workpiece having a complicated three-dimensional shaped recess is finished by electrolytic machining (see Japanese Patent Application Laid-Open 63-283818 and U.S. Pat. No. 4,800,006). In the system, since an electrode has a surface corresponding to the recessed surface of the workpiece, it is very difficult to center the workpiece with the electrode and a long time is consumed for the centering operation. Japanese Patent Application Publication 62-58850 discloses a centering method in which deflection of an electrode is automatically corrected. However, such a method can not be applied to the system employed with a pair of electrodes consisting of the workpiece having three-dimensional shaped recess and the electrode.\nThe object of the present invention is to provide a centering method in which the workpiece is a automatically centered with the electrode, thereby reducing operation time for centering."}
{"text": "Gene traps provide a general strategy to identify genes exhibiting discrete patterns of expression during development and differentiation. The basic design of gene trap vectors has been based on the introduction of a promoterless reporter gene, e.g., xcex2-galactosidase, into embryonic stem (ES) cells, which could only be expressed if the reporter gene has integrated in the right frame and orientation within a transcriptional unit. To improve the chances of expression, the reporter gene has been placed. downstream of a splice acceptor (SA) sequence allowing expression to occur when the reporter gene is integrated within an intron. Integration results in reporter gene expression that reflects the expression pattern of the endogenous gene or is influenced by nearby transcriptional regulatory elements.\nOne method of constructing trap vectors has been to use as a background a retroviral vector. Retroviruses integrate into the genome with no rearrangements of flanking sequences. This is not always the case when DNA is introduced by microinjection and perhaps other methods. An additional advantage of using a retroviral vector has been that the sites of proviral integration are often found close to hypersensitive sites (Vijaya et al., J. Virol. 60:683-692 (1986); Rohdewohld et al., J. Virol. 61:336-343 (1987)).\nIn one particular design the reporter gene is inserted in Reverse Orientation (with respect to retroviral transcription) downstream of a Slice Acceptor sequence, and the resultant mice are referred to as ROSA. Using this gene trap design several different reporter genes have been used in embryonic stem cells as a genetic screen to identify and mutate developmental genes in mice (Friedrich and Soriano, Genes and Development 5:1513-1523 (1991)). Using the ROSA gene trap design many gene trap transgenic lines have been derived that display various expression patterns in embryos at different developmental stages. Among the transgenic lines obtained have been certain promoters which appear to be ubiquitously expressed.\nGene targeting in murine embryonic stem cells has allowed the production of mice with specific gene deletions. This technique has been used in studies to try and determine the functional identification of gene products. Conventional gene knockout techniques have provided mice that inherit genetic deletions in all cell types in a regionally and temporally unrestricted manner which can lead to severe developmental defects and premature death of the knock-out animals. Several model systems have been developed which attempt to take advantage of the cell-type or tissue-type restricted expression of certain promoters operatively associated with a recombinase gene in combination with a gene of interest that has been flanked by recombinase recognition sequences. In this system the recombinase is expressed under the control of the cell-type or tissue-type specific promoter and when expressed results in the excision of the gene of interest. Because the temporal nature of cell-type and/or tissue-type specificity of the promoters is not known with any certainty results obtained with this system are suspect. Further, as these promoters are not active in all cells or in all tissues of an animal they are not-as useful for examining conditional mutations in genes.\nThe present invention provides methods and vector constructs which quite unexpectedly provide means to position a gene of interest under the control of an ubiquitously expressed promoter and a means to confirm the nature of the expression of the promoter sequence such that the study of these mutations in transgenic non-human organisms can be accomplished.\nThe present invention provides methods and vector constructs for the production of genetically engineered non-human animals which ubiquitously express a heterologous DNA segment. In one embodiment of the invention, the methods comprise transforming a pluripotent cell with a DNA construct comprising a heterologous DNA segment in which at least 100 base pairs is homologous with a DNA sequence of an ubiquitously expressed endogenous gene locus of the pluripotent cell, wherein the DNA construct becomes integrated into the gene locus by homologous recombination. Insertion of the DNA construct into the ubiquitously expressed gene locus places the DNA construct under the control of the promoter of the ubiquitously expressed gene locus and does not result in a lethal mutation.\nPluripotent cells which carry the heterologous gene inserted into and under the control of the ubiquitously expressed gene locus promoter are selected and introduced into a developing embryo, at, for example, the blastocyst or morula stage. The embryos are allowed to develop to term and offspring are selected which carry the heterologous DNA segment integrated into the ubiquitously expressed endogenous gene locus and under its promoter.\nIn a particularly preferred embodiment of the invention, the pluripotent cells are murine embryonic stem cells, zygotes or sperm cells, and the like. Also, in an alternative embodiment, a splice acceptor sequence is operatively associated with the heterologous DNA segment.\nThe methods of the present invention can be used to produce general deletor and general reporter (alternatively designated a Universal Conditional Reporter (UCR)), animal strains. In one representative embodiment of a general deletor animal, the heterologous DNA segment is a general deletor cassette comprising a gene encoding a recombinase. Optionally, a splice acceptor sequence can be associated with gene encoding the recombinase. The animal produced using the methods of the present invention ubiquitously express the recombinase in essentially all cells and all tissues throughout development. When these mice are crossed with an animal strain which has a gene of interest flanked by recombinase recognition sequences recognized by the recombinase expressed by the general deletor mouse, the gene of interest is excised from the chromosome. It is necessary for the two recombinase recognition sequences to be in the same orientation within a ubiquitously expressed endogenous gene locus. The heterologous DNA segment is positioned within the ubiquitously expressed endogenous promoter such that the reporter cassette expression is under the control of the promoter. A particularly preferred reporter is xcex2-galactosidase.\nRepresentative ubiquitously expressed endogenous gene loci which. can be used in the methods of the present invention include ROSA26, ROSA5, ROSA23, ROSA11, and G3BP (BT5). Other loci can be determined by gene traps or other well known methods.\nIn an alternative embodiment of the present invention, the general deletor cassette further comprises a positive selection cassette downstream of the heterologous DNA segment. This cassette can be used in the identification of pluripotent cells which have been integrated under control of a ubiquitous promoter. Representative positive selectable markers include neo, gpt and others.\nIn an alternative embodiment of the general reporter cassette, the DNA stuffer sequence can comprise a promoter operatively associated with a selectable marker. The promoter can be inducible if desired. One particularly preferred promoter is PGK. Particularly preferred recombinase recognition sequences are lox and frt which are recognized by Cre and Flp recombinase, respectively.\nThe present invention also provides a general targeting vector cassette which comprises at least 100 base pairs of a DNA sequence homologous with an ubiquitously expressed endogenous gene locus and optionally, a negative selection cassette.\nRepresentative gene loci which are ubiquitously expressed and can be used in the present invention include ROSA26, ROSA5, ROSA23, ROSA11 and G3BP (BTS). A particularly preferred negative selection marker is Diphtheria toxin. In a particularly preferred embodiment, 5 kb of the ROSA 26 locus is used to create a general targeting vector cassette.\nIn another embodiment of the present invention, a general deletor cassette is provided. The general deletor cassette comprises at least 100 base pairs of a DNA sequence homologous with an ubiquitously expressed endogenous gene locus and, optionally, a negative selection cassette. Positioned within the homologous DNA sequence is a gene encoding a recombinase. Preferred recombinase for use in the present invention are Cre and Flp. The recombinase can further be associated with an Internal Ribosome Entry Site (IRES) upstream of the recombinase gene. A particularly preferred IRES is derived from the Encephalomyocarditis virus.\nIn an alternative embodiment, a splice acceptor sequence is operatively associated with the recombinase gene. Downstream of the recombinase gene, the general deletor vector cassette can further include a positive selection cassette. In a preferred embodiment, the positive selection cassette comprises a promoter, such as PGK, and a positive selectable marker, such as neo or Herpes simplex virus tk. Particularly preferred ubiquitously expressed endogenous gene loci for use in the general deletor cassette are ROSA26, ROSA5, ROSA23, ROSA11, and G3BP (BT5).\nIn yet another embodiment of the present invention, a conditional deletor cassette is provided. The conditional deletor cassette comprises at least 100 base pairs of a DNA sequence homologous with a conditionally expressed endogenous gene locus, and optionally, a positive or negative selection cassette. Positioned within the homologous DNA sequence is a gene encoding a recombinase as described above. The recombinase can be further associated with an IRES upstream of the recombinase gene. A particularly preferred IRES is derived from the Encephalomyocarditis virus.\nIn an alternative embodiment, a splice acceptor sequence is operatively associated with the recombinase gene. Downstream of the recombinase gene, the general deletor vector cassette can further include a positive selection cassette. In a preferred embodiment, the positive selection cassette comprises a promoter, such as PGK, and a positive selectable marker, such as neo or Herpes simplex virus tk. A particularly preferred ubiquitously expressed endogenous gene loci for use in the conditional deletor cassette is EphA2.\nIn still another embodiment of the present invention, a general reporter vector cassette is provided. This cassette comprises at least 100 base pairs of a DNA sequence homologous with an ubiquitously expressed endogenous gene locus and, optionally, a negative selection cassette. Within the DNA sequence is inserted a DNA stuffer sequence flanked by two recombinase recognition sequences in the same orientation positioned upstream of a gene encoding a reporter. In one alternative embodiment of the present invention an IRES is positioned upstream of the gene encoding a reporter. In another alternative embodiment, a splice acceptor is operatively associated with the DNA stuffer sequence. The negative selection cassette is preferred to comprise a promoter operatively associated with a negatively selectable marker. A particularly preferred selectable marker is the Diphtheria toxin gene.\nPreferred recombinase recognition sequences include, but are not limited to lox and Frt which are recognition sequences for the recombinases Cre and Flp, respectively."}
{"text": "1. Technical Field of the Invention\nThe present invention relates to a system and method for processing incoming telephone calls, and particularly to a system and method for providing operator assistance by automatically accessing a number of customer directories.\n2. Background and Objects of the Invention\nThere is a general market requirement for companies and other business entities to provide operator assistance to facilitate the handling of incoming telephone calls. A number of companies, such as companies in the European market, provide operator assistance to their incoming callers both during normal business hours and non-business hours in which the company is closed or otherwise not conducting business. Providing operator assistance seven days a week, 24 hours a day is, for many companies, cost prohibitive.\nIn existing operator service provider systems, a distinct directory is maintained for each customer (company), with the directories stored in one or more electronic databases. An operator receiving an incoming call for a subscriber of the customer, such as an employee of a customer company, must manually switch to the appropriate customer directory or, in the event all customer directories are maintained in a single database, manually execute a directory search to identify the desired customer group.\nSuch existing systems, however, present a number of shortcomings. First, operator-initiated directory switching/searching is too slow. In addition, errors are frequently encountered in operator-initiated directory switching/searching due to operator confusion or mistake. Erroneous switching/searching results in the incoming caller being transferred to the incorrect subscriber or even the incorrect customer/company. Further, existing operator service systems require that customer directories be exclusively maintained by the service provider in order to allow for special customer directory functions. As a result, there is a need for a more efficient and error-free operator service.\nIt is an object of the present invention to provide a telephone call processing system and method for providing an operator service.\nIt is another object of the present invention to provide an efficient and reliable system and service for managing incoming telephone calls directed to a number of different customers.\nAnother object of the present invention is to provide substantially error-free searching for subscriber information within the appropriate customer directory."}
{"text": "1. Field of the Invention\nThe invention relates to pedestrian warning devices, particularly those that are incorporated into roadway corners or sidewalks for pedestrian use.\n2. Discussion of Background Information\nSidewalks are well known structures that provide a walking surface for pedestrians, and that are frequently used on the sides of roads, parking lots, and other similar areas where pedestrians and vehicles are both routinely present. Such sidewalks are typically made from a form of pavement and are typically anywhere from 4 to 8 inches above the road surface.\nTo make sidewalks safer, particularly for users with disabilities, ramps are used to create a sloped transition from the sidewalk to the road surface. These ramps and associated landings are subject to numerous government regulations, particularly under the Americans with Disabilities Act, that are intended to ensure that the surface is safe for use by pedestrians and disabled persons.\nIn general, these ramps or landings have a concrete base, similar to the sidewalk itself, that is covered with a tactile panel, such as the Universal Radius Tactile Warning Surface Product (U.S. Pat. No. 9,365,985 B2), which is a glass and carbon composite panel that is enhanced with fiberglass truncated domes. Traditionally, tactile panels are either cast-in-place or surface applied.\nAlthough cast-in-place tactile panels are easy to install in fresh concrete, they are difficult and costly to replace because once the concrete is set, the underlying substrate must be at least partially destroyed in order to remove and replace the tactile panel. One variation of cast-in-place tactile panels features spaced honeycomb-like lower walls. However, because of the nature of the honeycomb, air can become trapped between the lower walls, creating areas that lack support from the underlying substrate. This trapped air can result in a lack of support that can lead to fatigue, and cracking failure of the tactile panel due to repetitive and heavy loading.\nSurface applied tactile panels are tactile panels applied to a finished substrate. Surface applied tactile panels are typically mechanically fastened and adhered to the underlying substrate. In addition, in order to compensate for irregularities in the substrate and minimize water intrusion that would otherwise damage the tactile panel, the surface applied tactile panels require additional caulking around the perimeter. Because surface applied tactile panels are mechanically fastened, they are easier and cheaper to remove than cast-in-place tactile panels. Once the fasteners are removed, the panel is heated to break the adhesive bond then stripped away from the underlying substrate. As a result, the underlying substrate remains intact, reducing the cost and time associated with replacing the tactile panel.\nAlthough the surface applied tactile panel improves upon the cast-in-place tactile panel, both cast-in-place and surface applied tactile panels suffer from similar shortcomings. Both fail to remove the potential for air and water infiltration, which in turn compromises the structural integrity and functionality of the tactile panels.\nThese shortcomings are particularly problematic because the intersections of pedestrian ramps and sidewalks and street gutters are often locations where it is difficult to drain away standing water that remains after a rain or snow-melt event. Frequently there is insufficient pitch in these areas to provide positive drainage relief and puddles and ice may form resulting in hazards to pedestrians. If water is able to enter around the tactile panel at such a location or between the panel and the concrete slab, over time it will expand or contract, thus causing numerous problems to the long-term safety and stability of the ramp. To address these issues, ramp manufacturers go to great lengths to create water-tight or water-resistant tactile panels to prevent water from entering the ramp system. However, with a water tight ramp, pools of standing water may form at the base of the ramp during times of rain or melting snow or ice, causing puddles to form in the roadways and often back up the ramp, thus introducing a new hazard.\nPorous concrete, also known as porous pavement, is distinguishable from traditional or non-porous concrete due to its high porosity, which allows for stormwater to infiltrate back into the ground naturally by passing directly through the concrete. Porous concrete has the ability to reduce pavement runoff and prevent pooling in areas that lack adequate drainage. As a result, porous concrete slabs allow water, whether the water comes from rain, melting snow or ice, or general road cleaning activities, to drain through the porous concrete and infiltrate back into the ground. However, current tactile panels do not allow water flow, thereby negating the benefits of using porous concrete as an underlying substrate.\nWhat is needed, therefore, is a tactile warning device that is able to reduce the pooling or collection of water on or around the ramp and does not suffer structural degradation due to air and water infiltration."}
{"text": "Opiates are most commonly used as analgesics for treating severe chronic and acute pain. The opiate most frequently used at present for treating severe chronic or acute pain is morphine. An example of chronic pain is the pain which occurs in cancer. An example of acute pain is the pain which may occur after operations. The opiates used up until now for treating such pain are indeed highly effective but have a number of unpleasant and/or undesirable side effects, e.g. a short duration of activity, respiratory depression, nausea, constipation, diuresis and euphoria and they are also addictive.\nMorphine hydrazone derivatives with analgesic activity are known from EP-A-0 242 417. EP-A-0 577 847 discloses 6-N-substituted morphine derivatives with analgesic and diuretic activity. EP-A-0 632 041 discloses 6-nicotinoylaminomorphine derivatives having analgesic activity.\nA number of publications disclose preparations which attempt to avoid some of the known disadvantages of the opiates used up until now. EP-B-300 806 discloses the use of phospholipid vesicles for encapsulating opioid analgesics. EP-A-672 416, EP-A-647 448, EP-A-631 781 and WO 94/22431 disclose long-acting formulations of opiates in a hydrophobic matrix. All the formulations mentioned above have a duration of activity of from 12 to 24 hours. The disadvantages of these preparations are the delayed start of activity and the side effects, which still occur."}
{"text": "1. Field of the Invention\nThe present invention relates to locks, especially electronic locks having motor-driven bolts. More specifically, the invention relates to locks in which it is desired that the bolt, once extended, cannot be forcibly pushed in but can only be withdrawn into the lock with entry of a proper combination or other authorization. The invention also relates to locks in which it is desired to respond to certain types of physical attack by rendering the bolt incapable of being withdrawn. The invention further relates to locks in which various security enhancements are provided.\n2. Related Art\nNumerous conventional lock designs have been provided in which a bolt may be extended or withdrawn in response to entry of a combination of other authorization. However, some of the designs have not provided a \"dead bolt\" feature, which involves physical blocking of the extended bolt so that, after the bolt has been extended into its \"locked\" position, the lock resists externally applied pressure that attempts to force the bolt back into the lock case.\nAlso, it is envisioned that locks are physically attacked in may ways, including drilling into the lock case. It is desired that a lock not merely physically resist such attacks, but also respond appropriately to such attacks by ensuring that the bolt cannot be withdrawn during or after the attack. In other words, it is desirable to prolong or perpetuate the \"dead bolt\" state so that in the event of physical attack, it becomes even harder for a perpetrator to gain entry into the protected area. Many known locks do not prolong or perpetuate a \"dead bolt\" state after the lock has been physically attacked, and thus do not provide adequate additional protection in that scenario.\nFurther, many known lock systems that involve \"bolt works\" require two separate actions to extend the blocking member from the door into the doorjamb, and to re-extend the bolt from the lock case. This is not merely inconvenient, but presents an additional security risk should the individual neglect to perform the second action. Additionally, it is desirable in such systems to provide a \"bolt throw\" (extent of movement of the bolt) that is adjustable so as to easily adapt a single lock to a variety of installations and different types of bolt works.\nMoreover, many known lock systems possess minimal locking functions, and do not provide additional security enhancement features. Applicants have recognized that such security enhancements include detection and response to tampering with the keypad unit, remote enablement and disablement of the lock, detection and response to a user's attempting to open the lock while under duress, and the ability to store and later transmit a history of occurrences in the lock system.\nIt is to meet these and other goals that the present invention is directed. No known conventional lock is believed to have the features and advantages of the inventive locks that are described in the following specification."}
{"text": "Advances in electronic devices generally include reducing the size of the components that form integrated circuits. As semiconductor devices have become more highly integrated in recent years, circuit interconnections have become finer and distances between those circuit interconnections have become smaller. For example, in the case of photolithography, which can form interconnections smaller than 0.5 μm wide, surfaces on which pattern images are to be focused should be as flat as possible. Further, with smaller circuit components the value of each unit area of a semiconductor wafer becomes higher because the ability to use all of the wafer area for integrated circuit components improves. To properly form an integrated circuit that employs a much higher percentage of usable wafer area, contaminant particle counts on the semiconductor wafer surface must be reduced to very low levels. In order to clean a semiconductor wafer and remove unwanted particles, a process known as chemical mechanical polishing or chemical mechanical planarizing (hereinafter “CMP”) has become popular.\nThe polishing is typically accomplished using a polishing pad attached to a platen surface of a CMP machine. A workpiece to be polished, such as a semiconductor wafer, is arranged in a manner such that a surface to be polished faces the polishing pad. The polishing pad and workpiece are independently rotated while an abrasive liquid or slurry is supplied onto the polishing pad. The workpiece is then pressed against the polishing pad at a predetermined pressure and the surface of the workpiece is polished to a flat mirror finish.\nPolishing pads have a limited service life and become less effective over time. In turn, this results in inconsistent quality in the semiconductor wafers. As the polishing pad begins to lose its effectiveness, the semiconductor wafers polished later in the pad's life may be less smooth, possibly resulting in significant amounts of lost or unsaleable product. Thus, the polishing pads must be changed from time to time.\nCurrent methods of determining the wear of a polishing pad, and thus whether the polishing pad is in need of replacement, are unsatisfactory. Analysis of the pad typically requires visual inspection. However, the pad usually cannot be visually inspected without removing the polishing pad from the CMP machine. Once removed, the polishing pad typically cannot be placed back on the machine, even if visual inspection determines some usable life remains. As a result, statistical analysis, based on experimental determination of pad wear over time under various operating conditions, is often used to make decisions about pad replacement. This is unsatisfactory for the reason, among others, that statistical analysis represents an average based on the wear of previous pads and is unable to account for any one particular pad or variations in quality or aberrant pads of inferior quality.\nFurthermore, analysis of each polishing pad typically takes several hours or more and is often conducted by highly paid, highly skilled technicians whose time might be better used on other projects.\nAnother method of monitoring and analyzing polishing pad effectiveness is performing an analysis of the smoothness of the polished semiconductor wafer product. However, due to other steps required in wafer manufacture, this analysis is impractical until later in the manufacturing process. Thus, by the time defects in any one polishing pad are discovered and needed replacement is determined, the worn out pad may already have improperly polished many wafers remaining upstream of a suitable wafer analysis point, resulting in significant loss of product that must be discarded or recycled.\nConversely, using methods such as statistical analysis may also result in removing the polishing pad too soon, meaning that the full usable life of the polishing pad has not been obtained, resulting in the inefficient use of consumable materials. Thus, this also results in undesirable and uneconomic business practices.\nWhat is needed is a system that can provide more efficient analysis of CMP polishing pads. What is also needed is a system that can analyze CMP polishing pads without removing the pad from the CMP machine."}
{"text": "The present invention is in the field of pharmacotherapy. More specifically, it relates to the treatment of diseases and conditions affecting the nails of humans and other mammals.\nWhile many cutaneous disorders affect the skin itself, there also exist diseases and conditions which relate to the skin appendages, in particular to the nails. These are often difficult to treat due to the thick and dense nail plate largely composed of keratin and its poor uptake of therapeutic agents.\nNails are hardenings of the horny zone of the epidermis. They appear as sheet-like appendages covering the skin on the dorsal side of the terminal phalanges of fingers and toes. The horny zone of the nail is composed of hard alpha-keratin and has a distal, exposed part, or body, and a proximal, hidden portion, or root. The root is covered by a distal prolongation of the stratum corneum of the skin. This narrow fold is composed of soft-keratin and is termed the eponychium. Distal to the eponychium is the “half-moon,” or lunula, a part of the horny zone that is opaque to the underlying capillaries.\nDeep to the distal or free border of the nail, the horny zone of the fingertip is thickened and is frequently termed the hyponychium. The horny zone of the nail is attached to the underlying nail bed. The matrix, or proximal part of the bed, produces hard-keratin. Further distally, however, the bed may also generate nail substance. Moreover, the most superficial layer of the nail may be produced by the epithelium immediately dorsal to the root and proximal to the eponychium. The growth of the nail is affected by nutrition, hormones, and disease. Nail growth involves considerable protein synthesis, as a result of which nonspecific changes occur in the nails in response to various local and systemic disturbances.\nNails develop in the fetus as epidermal thickenings that undercut the skin to form folds from which the horny substance of the nail grows distally. In adult humans, it takes about 6 months for a fingernail to form, which corresponds to a growth rate of approx. 2-3 mm per month. Toenails tend to grow more slowly than fingernails.\nNails are about two magnitudes thicker than the stratum corneum of the skin. Hard alpha-keratin, the major constituent of the horny zone, is a fibrous structural protein characterised by a high content of cysteine which readily forms thermostable crosslinks via sulphide bridges. The water content of nails is rather low, and typically ranges from about 10 to 30%. Generally speaking, nails are very resistant to the permeation of molecules such as drug substances.\nThe main function of nails is to protect the sensitive tips of fingers and toes. In addition, fingernails serve in scratching. Toenails are also important for balance.\nOne of the most common diseases affecting the nail apparatus is fungal infection, also referred to as onychomycosis, a condition affecting about 20% of the adult population in the USA. The incidence is increasing worldwide. In fact, about 30% of all superficial fungal infections affect the nail. Infection may be due to dermatophyte (ringworm, tinea unguium), yeast, or other non-dermatophyte (mould) species. In paronychia, chronic infection of the nail fold is most often caused by Candida species, but bacterial infection with Gram negative species such as Pseudomonas may coexist. Acute paronychia (whitlow) due to staphylococcal infection may also occur, and the presence of these bacterial infections will influence management. Invasion of the nail plate by Candida species may occur in the presence of paronychia, immune deficiency states (including chronic mucocutaneous candidiasis), Raynaud's disease, or endocrine disorders.\nAnother common condition is nail psoriasis, one of the possible manifestations of psoriasis, which produces a variety of changes in the appearance of finger and toe nails. These changes include discolouring under the nail plate, pitting of the nails, lines going across the nails, thickening of the skin under the nail, and the loosening (onycholysis) and crumbling of the nail.\nMost people who have psoriasis of the nails also have skin psoriasis (cutaneous psoriasis). Only 5% of people with psoriasis of the nails do not have skin psoriasis. In people who have skin psoriasis, 10%-55% have psoriasis of the nails (also called psoriatic nail disease). About 10%-20% of people who have skin psoriasis also have psoriatic arthritis, a specific condition in which people have symptoms of both arthritis and psoriasis. Of people with psoriatic arthritis, 53%-86% have affected nails, often with pitting.\nAnother condition of the nail is onychia, which is an inflammation of the nail folds surrounding tissue of the nail plate with formation of pus and shedding of the nail. One of the causes of onychia is a posttraumatic bacterial infection.\nOnycholysis refers to a loosening of the exposed portion of the nail from the nail bed, usually beginning at the free edge and continuing to the lunula. It may be related to numerous possible causes including infections and allergic reactions.\nA number of other nail conditions such as onychodystrophy, onychomadesis, and onychoptosis often represent adverse reactions to drugs, for example antibiotics or anticancer agents.\nInterestingly, conditions of the nails are rather often treated systemically, which itself indicates how difficult it is to achieve a therapeutic drug concentration in the nails by local administration. For example, onychomycosis is treated with oral terbinafine 250 mg daily for 3-6 months, or itraconazole 200 mg daily for 3-6 months. Even the rather conventional therapy with oral griseofulvin 10 mg/kg/daily (500 mg twice daily) for 6-18 months is still being recommended today. In particular proximal nail disease or severe nail bed involvement are being considered as indications for systemic rather than local treatment.\nSome topical preparations of known antifungal agents for treatment of onychomycosis exist, such as nail lacquers comprising ciclopirox, amorolfine, or butenafine. While there is some evidence of efficacy, it is believed that successful topical antifungal therapy requires treatment over very long periods, such as a year or even more. Some experts recommend the combination of topical and systemic treatment for better efficacy. Without adequate treatment and patient compliance, the infection will not disappear.\nAlso for nail psoriasis, the most effective treatment options appear to be systemic rather than topical. Recent reports suggest that some injectable biological medicines such as infliximab and etanercept may be highly effective. However, these treatments are associated with substantial risks of adverse effects, and they are very expensive.\nTopical treatments in their currently available formulations are considered much less effective. Solutions containing vitamin D derivatives such as calcipotriol may be applied twice daily to the nail folds. Topical high-potency corticosteroid solutions or ointments are another option. Some experts also recommend 5-fluorouracil cream applied to the matrix for 6 months to improve pitting and subungual hyperkeratosis. Antifungal treatment is indicated whenever a secondary fungal infection is present, which is very frequent in cases of nail psoriasis.\nThere have been various efforts to make topical treatment of nail diseases more successful, mostly focussing on an improved permeation of the drug substance into the nail matrix. Often, co-treatment with urea cream is recommended in order to soften the nail keratin and make it more permeable. While some evidence for an improved antifungal therapy in combination with urea exists, it still appears that systemic therapy with all its risks and adverse effects is considerably more effective.\nOther permeation enhancers have been proposed for enhancing trans-nail delivery of antifungal agents, such as in U.S. Pat. No. 6,042,845, U.S. Pat. No. 6,159,977, U.S. Pat. No. 6,224,887 and U.S. Pat. No. 6,391,879. Most of the suggested compounds are however not approved for pharmaceutical products and exhibit unknown health risks.\nU.S. Pat. No. 5,326,566 describes a composition of a pharmacological agent in combination with dibutyl adipate, or a mixture of dibutyl adipate and isopropyl myristate, which could enhance the penetration through keratin. However, substantial problems may arise when the penetration enhancers are incompatible with a particular drug substance, leading to drug instability and degradation into potentially harmful degradants.\nUS 2005/0079210 A1 proposes the use of liposomes the epicutaneous administration of drugs and cosmetically useful agents. However, liposomes are difficult to manufacture cost-effectively and in a reproducible manner.\nOther delivery strategies for topical medications require the occlusion of the nail after administration, which is however perceived as inconvenient by many patients.\nThere clearly remains a need for pharmaceutical formulations and vehicles which allow the effective treatment of nail disorders by topical administration. It is therefore an object of the present invention to provide such improved compositions which overcome one or more disadvantages of known compositions. In particular, it is an object of the invention to provide compositions of topically active drug substances which can effectively penetrate the dense nail plate and which are convenient to use. Further objects of the invention will become clear on the basis of the description of the invention below, including the examples, and of the patent claims."}
{"text": "In a network device, a scheduler is commonly used to schedule the processing of data units (e.g. packets, cells, frames, etc.) associated with multiple traffic flows. Typically, a weighted round-robin (WRR) scheduler may process data units according to WRR weights associated with queues. Each WRR weight may indicate a level of priority, in which the larger the WRR weight, the higher the priority. For example, consider three queues, namely, queue 0, queue 1, and queue 2, which are assigned respective WRR weights of 10, 20, and 30. Given this configuration, the WRR scheduler services queue 0 for 10 timeslots, and queue 1 may have to wait those 10 timeslots before the WRR scheduler services queue 1. Additionally, queue 2 may have to wait 30 timeslots (i.e., 10 timeslots for queue 0 and 20 timeslots for queue 1) before the WRR scheduler services queue 2. As a result, the WRR scheduler may cause periods of burstiness to the traffic flow associated with, for example, queue 2, which may increase the amount of buffering and/or may cause a higher latency to occur with respect to queue 0 and/or queue 1."}
{"text": "The present invention generally relates to the inverting of devices, and particularly to the simultaneous inverting of a plurality of devices.\nDuring the manufacture of devices such as semiconductors and similar electronic devices, such devices must be tested for operability, efficiency, and other conditions. Thus, a need exists for apparatus which mechanically handles such devices to reduce the cost of such testing. Additionally, such mechanical handling also includes the need to perform tests or other functions on both sides of the device.\nWhen it is necessary to remove the device from the carrier or tray, inverting apparatus such as of the type disclosed in U.S. Pat. No. 5,374,158 has proven to be very effective. However, especially when it is not necessary to remove the devices from the carrier such as for mechanical vision testing, branding, or the like, inverting the devices one at a time is very time consuming and inefficient.\nThus, a need exists for simultaneously placing a plurality of devices contained in a first carrier in an inverted orientation. Past attempts to solve this need have not proven to be effective. Specifically, the positioning of the devices in the carrier is required for referencing the results of the various tests and other processing operations to individual devices. Because of warping of the carrier and for other reasons, devices tended to move in the carrier while they were being inverted so that damage to the devices such as to the leads thereof during the inverting process or during later handling operations and/or so that reference positioning was lost in the carrier. Additionally, devices also tended to fall from the carrier which is clearly undesirable due to the loss of the devices themselves but also because of the potential the devices could fall upon and short out electrical equipment in the vicinity."}
{"text": "Short range, license-free, wireless applications operating in the Industrial Scientific Medical (ISM) bands (e.g. door openers, tire pressure monitoring systems (TPMS), wireless mice and wireless local area networks (WLANs)) have become popular in recent years. Such applications typically comprise one or more transmitting devices, which transmit a radio signal to one or more receiving devices. On receipt of a signal from the transmitting device(s), the receiving device(s) take an appropriate action.\nTransmitting device(s) are typically mobile, handheld devices, which comprise a transmitter (for generating signals), an antenna (for transmitting the signals) and a battery (as an internal energy source for such transmissions). However, the batteries are usually quite small. Thus, power consumption is a critical factor in the design of such transmitting device(s).\nThe antenna in a transmitting device is a load on the device's transmitter. Maximally efficient power transfer between a transmitter and an antenna is only achieved when the transmitter impedance is the complex conjugate of the antenna impedance. However, the sensitivity of antenna impedance to surrounding conditions is making the task of matching the impedances of antennas and transmitters quite difficult. Attempts to solve this problem have been ongoing for some time.\nIntegration Associates (Automatic Antenna Tuning RF Transmitter IC Applying High Q Antenna—White Paper Ver. 1.0 (IA ISM-WP1), Integration Associates Inc., 2004) and Micrel (MICRF103 QwikRadio (trademark) ASK Transmitter FINAL, June 2002) have proposed a solution in which a single capacitor adjusts the reactive part of antenna impedance. However, any reduction in impedance mismatch achieved by this method would be limited, as it does not take into account the resistive components of the antenna impedance and transmitter impedance. Moreover, the Integration Associates and Micrel solutions are only applicable to source current transmitters, as the measurements used to match the impedances are phase differences between the base and collector of a transistor in a switched current source configuration.\nU.S. Patent Application US2005219132 describes a system in which two capacitors of a PI matching network are tuned. Impedance matching is achieved by measuring the voltage and the phase difference between the input and output of the matching network. However, the system described in US2005219132 requires A/D convertors and a phase comparator and would be difficult to implement at the frequency of 868 MHz typically employed in automotive applications.\nU.S. Patent Application US2005003771 describes a circuit for automatically tuning a resonant circuit in a transceiver. The circuit adjusts a single capacitor connected between the two outputs of a differential amplifier, wherein the adjustments are made on the basis of the phase difference between the input and the output of the amplifier. However, in a similar manner to the Integration Associates and Micrel solutions, the circuit in US2005003771 only provides reactive impedance matching. Furthermore, the circuit does not work with Class C amplifiers.\nMore generally, by using only a single adjustable capacitor, prior art solutions available at frequencies of 434-868 MHz have limited ability to match the impedances of a transmitter and antenna. In particular, whilst the prior art solutions provide reactive impedance cancellation, they do not provide conjugate matching."}
{"text": "This invention relates to treating seeds. More specifically, this invention relates to a device that provides chemicals to seed for treatment.\nFor years companies have attempted to improve agricultural seeds by adding chemicals to the seeds in order to enhance yield, growth, and the like. Specifically, in addition to utilizing fertilizers and herbicides in the field seeds can be chemically treated before planting in order to improve crops yields, plant health, and the like.\nCurrently in the art a machine is provided that receives seeds within a chamber where chemicals are injected within the chamber for treatment of the seeds. This is typically done by the utilization of a syringe or other flow device that rapidly places the chemical fluid into the chamber. Presently, the amount of fluid used within the chamber is determined by the weight of the seed that is inputted into the chamber.\nWhile this type of device is effective at introducing the seeds to the chemicals problems still remain. In particular, the chemicals utilized in treating the seeds are typically very expensive and a large amount of the chemicals are wasted during the treatment process. In particular, because a syringe is used and because the chemicals are added based on the weight of the seeds excess chemicals are pumped into the seed treating chamber thus creating waste and increased expense that is undesired. Thus, a need in the art exists for an improved device or method for treating seeds with chemicals.\nTherefore, a principal object of the present invention is to provide a seed treating device that reduces the amount of chemicals needed to treat seeds.\nYet another object of the present invention is to provide a seed treating device that is efficient and eliminates waste.\nThese and other objects, advantages and features will become apparent from the specification and claims."}
{"text": "1. Field of the Invention\nThis invention relates to a car power source apparatus provided with heaters to increase the temperature of the driving battery that powers the electric drive motor of a vehicle, such as an automobile, when temperatures are low.\n2. Background Art\nWhen temperatures become extremely low, battery output decreases and sufficient performance cannot be obtained. Consequently, when a hybrid car is used in a frigid location and the temperature of the driving battery becomes low, sufficient battery performance cannot be output when beginning to move or accelerate the car. A hybrid car is designed to normally run on both an engine and an electric motor. If performance of the driving battery decreases and power from the motor is insufficient, the hybrid car cannot operate normally. This drawback can be resolved by warming-up (heating) the battery. To realize this, the present applicant developed a power source apparatus which warms-up rechargeable batteries via heaters (see Japanese Laid-Open Patent Publication 2003-223938).\nThe power source apparatus disclosed in this patent publication is provided with a driving battery comprising a plurality of rechargeable batteries, and a heating plate, with heaters mounted on it, to warm-up (heat) the rechargeable batteries of the driving battery. The heating plate is disposed in close proximity to the rechargeable batteries to allow efficient heating of those rechargeable batteries. Electric current is passed through the heating plate heaters of this power source apparatus to heat the heaters via Joule heating. Many heaters are disposed in close proximity to many rechargeable batteries to heat those batteries."}
{"text": "1. Field of the Invention\nThe present general inventive concept relates to a computer system and a program restoring method thereof, and more particularly to a computer system storing with a restoration program and a restoring method using the restoration program.\n2. Description of the Related Art\nIn general, a computer system performs various functions through hardware such as a central processing unit (CPU), and software such as an operating system, a device driver, and an application. In the case of the software, update is often periodically provided by an external server or the like with development of its technology. However, the updated software may have a problem with its operation due to various factors, for example, if it is infected with a virus while using a computer system, deleted by mistake, modified, or etc. At this time, the updated software having the problem has to be restored or recovered, but a restoration solution initially provided by a manufacturer of a computer system provides only a previous version of software. Thus, a user has to access the external server or the like after restoring the software, and then update the restored software with the latest version."}
{"text": "The present invention is related to a co-pending application filed concurrently herewith and entitled xe2x80x9cScalable Internet Engine,xe2x80x9d a copy of which is attached hereto and the disclosure of which is hereby incorporated by reference.\nThe present invention relates generally to the field of data processing business practices. More specifically, the present invention relates to a method and system for providing dynamic management of hosted services across disparate customer accounts and/or geographically distinct sites.\nThe explosive growth of the Internet has been driven to a large extent by the emergence of commercial service providers and hosting facilities, such as Internet Service Providers (ISPs), Application Service Providers (ASPs), Independent Software Vendors (ISVs), Enterprise Solution Providers (ESPs), Managed Service Providers (MSPs) and the like. Although there is no clear definition of the precise set of services provided by each of these businesses, generally these service providers and hosting facilities provide services tailored to meet some, most or all of a customer\"\"s needs with respect to application hosting, site development, e-commerce management and server deployment in exchange for payment of setup charges and periodic fees. In the context of server deployment, for example, the fees are customarily based on the particular hardware and software configurations that a customer will specify for hosting the customer\"\"s application or website. For purposes of this invention, the term xe2x80x9chosted servicesxe2x80x9d is intended to encompass the various types of these services provided by this spectrum of service providers and hosting facilities. For convenience, this group of service providers and hosting facilities shall be referred to collectively as Hosted Service Providers (HSPs).\nCommercial HSPs provide users with access to hosted applications on the Internet in the same way that telephone companies provide customers with connections to their intended caller through the international telephone network. The computer equipment that HSPs use to host the applications and services they provide is commonly referred to as a server. In its simplest form, a server can be a personal computer that is connected to the Internet through a network interface and that runs specific software designed to service the requests made by customers or clients of that server. For all of the various delivery models that can be used by HSPs to provide hosted services, most HSPs will use a collection of servers that are connected to an internal network in what is commonly referred to as a xe2x80x9cserver farmxe2x80x9d, with each server performing unique tasks or the group of servers sharing the load of multiple tasks, such as mail server, web server, access server, accounting and management server. In the context of hosting websites, for example, customers with smaller websites are often aggregated onto and supported by a single web server. Larger websites, however, are commonly hosted on dedicated web servers that provide services solely for that site. For general background on the Internet and HSPs, refer to Geoff Huston, ISP Survival Guide: Strategies For Running A Competitive ISP, (1999).\nAs the demand for Internet services has increased, there has been a need for ever-larger capacity to meet this demand. One solution has been to utilize more powerful computer systems as servers. Large mainframe and midsize computer systems have been used as servers to service large websites and corporate networks. Most HSPs tend not to utilize these larger computer systems because of the expense, complexity, and lack of flexibility of such systems. Instead, HSPs have preferred to utilize server farms consisting of large numbers of individual personal computer servers wired to a common Internet connection or bank of modems and sometimes accessing a common set of disk drives. When an HSP adds a new hosted service customer, for example, one or more personal computer servers are manually added to the HSP server farm and loaded with the appropriate software and data (e.g., web content) for that customer. In this way, the HSP deploys only that level of hardware required to support its current customer level. Equally as important, the HSP can charge its customers an upfront setup fee that covers a significant portion of the cost of this hardware. By utilizing this approach, the HSP does not have to spend money in advance for large computer systems with idle capacity that will not generate immediate revenue for the HSP. The server farm solution also affords an easier solution to the problem of maintaining security and data integrity across different customers than if those customers were all being serviced from a single larger mainframe computer. If all of the servers for a customer are loaded only with the software for that customer and are connected only to the data for that customer, security of that customer\"\"s information is insured by physical isolation.\nFor HSPs, numerous software billing packages are available to account and charge for these metered services, such as XaCCT from rens.com and HSP Power from inovaware.com. Other software programs have been developed to aid in the management of HSP networks, such as IP Magic from lightspeedsystems.com, Internet Services Management from resonate.com and MAMBA from luminate.com. The management and operation of an HSP has also been the subject of articles and seminars, such as Hursti, Jani, xe2x80x9cManagement of the Access Network and Service Provisioning,xe2x80x9d Seminar in Internetworking, Apr. 19, 1999. An example of of a typical HSP offering various configurations of hardware, software, maintenance and support for providing commercial levels of Internet access and website hosting at a monthly rate can be found at rackspace.com.\nUp to now, there have been two approaches with respect to the way in which HSPs built their server farms. One approach is to use a homogenous group of personal computer systems (hardware and software) supplied from a single manufacturer. The other approach is to use personal computer systems supplied from a number of different manufacturers. The homogeneous approach affords the HSP advantages in terms of only having to support a single server platform, but at the same time it restricts the HSP to this single server platform. The heterogeneous approach using systems supplied from different manufacturers is more flexible and affords the HSP the advantage of utilizing the most appropriate server hardware and software platform for a given customer or task, but this flexibility comes at the cost of increased complexity and support challenges associated with maintaining multiple server platforms.\nRegardless of which approach is used to populate a server farm, the actual physical management of such server farms remains generally the same. When a customer wants to increase or decrease the amount of services being provided for their account, the HSP will manually add or remove a server to or from that portion of the HSP server farm that is directly cabled to the data storage and network interconnect of that client\"\"s website. In the case where services are to be added, the typical process would be some variation of the following: (a) an order to change service level is received from a hosted service customer, (b) the HSP obtains new server hardware to meet the requested change, (c) personnel for the HSP physically install the new server hardware at the site where the server farm is located, (d) cabling for the new server hardware is added to the data storage and network connections for that site, (e) software for the server hardware is loaded onto the server and personnel for the HSP go through a series of initialization steps to configure the software specifically to the requirements of this customer account, and (f) the newly installed and fully configured server joins the existing administrative group of servers providing hosted service for the customer\"\"s account. In either case, each server farm is assigned to a specific customer and must be configured to meet the maximum projected demand for services from that customer account.\nOriginally, it was necessary to reboot or restart some or all of the existing servers in an administrative group for a given customer account in order to allow the last step of this process to be completed because pointers and tables in the existing servers would need to be manually updated to reflect the addition of a new server to the administrative group. This requirement dictated that changes in server hardware could only happen periodically in well-defined service windows, such as late on a Sunday night. More recently, software, such as Microsoft Windows 2000, Microsoft Cluster Server, Oracle Parallel Server, Windows Network Load Balancing Service (NLB), and similar programs have been developed and extended to automatically allow a new server to join an existing administrative group at any time rather than in these well-defined windows.\nAn example of how a new server can automatically join an existing administrative group is described in U.S. Pat. No. 5,951,694. In this patent, all of the servers in an administrative group are represented in a mapping table maintained by a gateway server. The mapping table identifies different service groups for the administrative group, such as mail service group, database service group, access server group, etc. The gateway server routes requests for the administrative group to the appropriate service group based on the mapping table. A new server may be added to one of the service groups by loading the appropriate software component on that server, after which the gateway server will recognize the new server and add it to the mapping table and bring the new server up to speed with the rest of the servers in that service group using a transaction log maintained for each service group. Alternatively, if one service group is experiencing a heavy workload and another service group is lightly loaded, it is possible to switch a server from one service group to another. The patent describes a software routine executing on a dedicated administrative server that uses a load balancing scheme to modify the mapping table to insure that requests for that administrative group are more evenly balanced among the various service groups that make up the administrative group.\nNumerous patents have described techniques for workload balancing among servers in a single cluster or administrative groups. U.S. Pat. No. 6,006,529 describes software clustering that includes security and heartbeat arrangement under control of a master server, where all of the cluster members are assigned a common IP address and load balancing is preformed within that cluster. U.S. Pat. Nos. 5,537,542, 5,948,065 and 5,974,462 describe various workload-balancing arrangements for a multi-system computer processing system having a shared data space. The distribution of work among servers can also be accomplished by interposing an intermediary system between the clients and servers. U.S. Pat. No. 6,097,882 describes a replicator system interposed between clients and servers to transparently redirect IP packets between the two based on server availability and workload.\nVarious techniques have also been used to coordinate the operation of multiple computers or servers in a single cluster. U.S. Pat. No. 6,014,669 describes cluster operation of multiple servers in a single cluster by using a lock-step distributed configuration file. U.S. Pat. No. 6,088,727 describes cluster control in a shared data space multi-computer environment. Other patents have described how a single image of the input/output space can be used to coordinate multiple computers. U.S. Pat. No. 5,832,222 describes how a single image of the input/output space can be used to coordinate geographically dispersed computer systems. U.S. Pat. No. 6,067,545 describes a distributed file system with shared metadata management, replicated configuration database and domain load balancing, that allows for servers to fall into and out of a single domain under control of the configuration database.\nWhile these approaches have improved the management of servers within administrative groups, domains or shared data spaces, there is no capability to extend these techniques beyond the group of servers defined for and linked to a common operating system or common shared data space. Generally, this limitation has not been considered a problem because all of these approaches are directed to larger enterprise computing systems that are managed and implemented within the computer network of a single company. Even though these approaches can be put into use by an HSP to manage the servers assigned to a particular account for a given client or customer, none of these approaches allow an HSP to manage a set of servers providing hosted services to multiple accounts for different clients or customers.\nSystems for managing the operation of larger enterprise computing systems also have been developed, such as OpenView from Hewlett-Packard, Unicenter TNG from Computer Associates, Tivoli from IBM, Mamba from Luminate, and Patrol from BMC Software, Inc. Generally, these systems are focused on inventory management and software deployment control issues encountered with very large numbers of computers operating within a single company or organization. Some of these operation management systems include performance monitoring solutions that query the performance of servers within the organization over the network to determine the need for additional resources or load redistribution. A similar over-the-network approach is also used to provide centralized reporting and management features. A good example of this type of operation management system that is intended to be used by HSPs is the Tivoli Service Delivery Management platform that consists of a user administration module, a software distribution module, an inventory module, an enterprise console, a security module, an enterprise manager module that provides a customizable view of all of the components in a network once they are added to the network, and a workload scheduler that allows workload to be balanced among servers sharing a common data space. All of these modules operate using an over-the-network communication scheme involving agents on the various nodes in the network that collect and report status and incident information to the other modules. Once the hardware components for a new node are physically added to the network, the various modules of the Tivoli Service Delivery Management platform can take over and manage those components on a more automatic basis. However, the process of physically adding hardware for a new node into the network remains essentially a manual process that is accomplished in the same manner as previously described.\nIn terms of managing the physical hardware that makes up the computer system, various approaches have been developed to automatically compensate for the failure of a hardware component within a computer network. U.S. Pat. No. 5,615,329 describes a typical example of a redundant hardware arrangement that implements remote data shadowing using dedicated separate primary and secondary computer systems where the secondary computer system takes over for the primary computer system in the event of a failure of the primary computer system. The problem with these types of mirroring or shadowing arrangements is that they can be expensive and wasteful, particularly where the secondary computer system is idled in a standby mode waiting for a failure of the primary computer system. U.S. Pat. No. 5,696,895 describes one solution to this problem in which a series of servers each run their own tasks, but each is also assigned to act as a backup to one of the other servers in the event that server has a failure. This arrangement allows the tasks being performed by both servers to continue on the backup server, although performance will be degraded. Other examples of this type of solution include the Epoch Point of Distribution (POD) server design and the USI Complex Web Service. The hardware components used to provide these services are predefined computing pods that include load-balancing software, which can also compensate for the failure of a hardware component within an administrative group. Even with the use of such predefined computing pods, the physical preparation and installation of such pods into an administrative group can take up to a week to accomplish.\nAll of these solutions can work to automatically manage and balance workloads and route around hardware failures within an administrative group based on an existing hardware computing capacity; however, few solutions have been developed that allow for the automatic deployment of additional hardware resources to an administrative group. If the potential need for additional hardware resources within an administrative group is known in advance, the most common solution is to preconfigure the hardware resources for an administrative group based on the highest predicted need for resources for that group. While this solution allows the administrative group to respond appropriately during times of peak demand, the extra hardware resources allocated to meet this peak demand are underutilized at most other times. As a result, the cost of providing hosted services for the administrative group is increased due to the underutilization of hardware resources for this group.\nOne solution to the need for additional hosted services is the Internet Shock Absorber (ISA) service offered by Cable and Wireless. The ISA service distributes a customer\"\"s static Web content to one or more caching servers located at various Points of Presence (POPs) on the Cable and Wireless Internet backbone. Requests for this static Web content can be directed to be caching servers and the various POP locations to offload this function from the servers in the administrative group providing hosted service for that customer. The caching of static Web content, however, is something that occurs naturally as part of the distribution of information over the Internet. Where a large number of users are requesting static information from a given the IP address, it is common to cache this information at multiple locations on the Internet. In essence, the ISA service allows a customer to proactively initiate the caching of static Web content on the Internet. While this solution has the potential to improve performance for delivery of static Web content, this solution is not applicable to the numerous other types of hosted services that involve interactive or dynamic information content.\nAlthough significant enhancements have been made to the way that HSPs are managed, and although many programs and tools have been developed to aid in the operation of HSP networks, the basic techniques used by HSPs to create and maintain the physical resources of a server farm have changed very little. It would be desirable to provide a more efficient way of operating an HSP that could improve on the way in which physical resources of the server farm are managed.\nThe present invention is a method and system for operating a hosted service provider for the Internet in such a way as to provide dynamic management of hosted services across disparate customer accounts and/or geographically distinct sites. For each of a plurality of customer accounts, a plurality of individual servers are allocated to a common administrative group defined for that customer account. Each administrative group is configured to access software and data unique to that customer account for providing hosted services to the Internet for that customer account. The system automatically monitors the performance and health of the servers in each administrative group. At least one server from a first administrative group is automatically and dynamically reallocated to a second administrative group in response to the automatic monitoring. The automatic and dynamic reallocation of servers is accomplished by setting initialization pointers for the reallocated servers to access software and data unique to the customer account for the second administrative group, and then reinitializing the reallocated servers such that they join the second administrative group when restarted. Preferably, the performance and health of the servers in each administrative group are monitored over a separate out-of-band communication channel dedicated to interconnecting the servers across administrative groups. Each administrative group includes a local decision software program that communicates with a master decision software program that determines when and how to dynamically reallocate servers to different administrative groups in response to usage demands, available resources and service level agreements with each customer account.\nIn one embodiment, a system for providing the dynamic management of hosted services for multiple customer accounts includes at least five servers operably connected to an intranet. Each server includes host management circuitry providing a communication channel with at least one of the other servers that is separate from this intranet. At least four of the servers execute a local decision software program that monitors the server and communicates status information across the communication channel. At least two of the servers are allocated to a first administrative group for a first customer account and configured to access software and data unique to this first customer account, such that hosted services are provided via the Internet for this customer account. At least two of the other servers are allocated to a second administrative group for a second customer account and configured to access software and data unique to this second customer account, such that hosted services are provided via the Internet for this customer account. Preferably, at least one of the servers executes a master decision software program that collects status information from the other servers and dynamically reallocates at least one server from the first administrative group to the second administrative group in response to at least the status information.\nUnlike existing load balancing systems that are limited to working within the context of a single customer account or that require large and expensive computer systems and common operating systems or shared data spaces, the present invention is capable of dynamically reallocating servers across multiple disparate customer accounts to provide hosted services with a more economical and flexible server farm arrangement. The ability of the present invention to support multiple administrative groups for multiple customers allows for an intelligent and dynamic allocation of server resources among different customer accounts."}
{"text": "The present disclosure relates to providing and enforcing security procedures in computer systems in general and more specifically relates to providing and enforcing security procedures in mobile information systems. In a mobile information system, subscribers may register with a wireless service provider to receive various types of content from the service provider on their mobile devices. The subscriber's mobile device may include a resident interactive multimedia application environment that includes capabilities for displaying graphics, video, animation, audio, and the like. Examples of such interactive multimedia application environments are the different versions of the Adobe® Flash®-based platform. Content provided to mobile devices equipped with such application environments is sometimes delivered in executable file formats such as the precisely described SWF (small web format) binary vector graphics format. SWF provides a compact, TAG-based, easily extendible format that supports streaming, bitmap and vector graphics, and scripting."}
{"text": "In manufacturing integrated circuit elements, there has been progress in miniaturization of processing size in lithography processes with a view to obtaining integrated circuits with a high degree of integration. In these lithography processes, a resist pattern is formed by coating a photoresist composition on a workpiece layer, followed by exposure and development, and the resist pattern thus formed is transferred to a workpiece film such as a wiring layer, a dielectric layer or the like. Conventionally, an exposed region of the workpiece film, being the region exposed in the resist pattern, was removed by dry etching. However, film thickness of the resist layer (i.e., the resist pattern) was decreased by shortening the wavelength of the light source for exposure or the like, accompanied by miniaturization of the processing size. Therefore, sufficient dry etching resistance cannot be acquired, thereby making the processing of the workpiece film with high accuracy difficult to achieve. Also, due to a problem of reflection of the exposed light on the workpiece film, it is difficult to form a resist pattern in a favorable shape.\nHence, for transferring a resist pattern to a workpiece film with good accuracy, insertion of an antireflective film (hard mask) between the workpiece film and the photo resist layer has been studied. As a characteristic of this antireflective film the etching rate is required to have a large difference from that of the resist pattern (see, for example, Patent Documents 1 to 3).\nPatent Document 1: Japanese Unexamined Patent Application Publication No. 2004-310019\nPatent Document 2: Japanese Unexamined Patent Application Publication No. 2005-015779\nPatent Document 3: Japanese Unexamined Patent Application Publication No. 2005-018054"}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, and more particularly to a multichip module (referred to as MCM hereinafter) semiconductor device formed by mounting, within the same package, a plurality of electronic components including semiconductor integrated circuit chips (referred to as IC chips hereinafter) into which desired functions are incorporated.\n2. Description of the Prior Art\nIn order to achieve further downsizing, weight reduction, thin forming and high performance for various kinds of apparatus employing a semiconductor device, high density packaging of various kinds of electronic component including IC chips has been investigated. As one of powerful means for achieving the objective there have been proposed various kinds of MCM semiconductor devices having a plurality of IC chips in the same package.\nFor example, a thin MCM package realizing a multichip semiconductor device with excellent heat dissipation has been proposed in Japanese Patent Applications Laid Open, No. Hei 8-250652, and an MCM semiconductor device aiming at low cost by employing a package that can be used universally even for different kinds or layouts of the IC chips to be mounted is proposed in Japanese Patent Applications Laid Open, No. Hei 9-181256.\nFIGS. 7A and 7B show schematic sectional views of examples of semiconductor devices using an MCM package disclosed in Japanese Patent Applications Laid Open, No. Hei 8-250652.\nFIG. 7A shows a diagrammatic drawing of a first example of MCM package 710. A printed wiring (PW) board 711 having a large number of level sections consists of a lower level section 712, an intermediate level section 713 and an upper level section 714. In this example, the lower level section 712 is continuous, whereas the intermediate and upper level sections respectively have through openings by means of which stepwise openings 715 are formed, and the stepwise openings form a cavity section 716 together with the lower level section 712.\nAn MCM tile is a silicon-on-silicon MCM tile 717 consisting of a silicon substrate 718 and silicon chips 719 and 720, and is situated in the cavity. The silicon substrate has a form in which it is placed in the cavity section on the surface of the lower level section of the PW board. Each of wire bond fingers 721 is interconnected to a contact pad 723 on the intermediate level section of the PW board via a wire 722.\nIn turn, each of these pads is interconnected to a part of another level section of the PW board, for example, a contact 725 via a through hole 724, thereby is interconnected to a solder bump 726 on the bottom face of the lower level section, and is interconnected as needed to another chip or an electronic device such as that represented by a symbol 727 or 728 on the surface of the upper level section of the PW board. Here, the MCM tile placed on the surface of the lower level section is located completely within the cavity section 716, and the top faces of the chips are at heights lower than the top face of the upper level section of the PW board.\nAn encapsulation sealing material 729 having high adaptability such as silicone gel is filled in the cavity section 716. The encapsulation sealing material 729 encapsulates the interconnecting sections between the chips and the silicon substrate, and the wire bond fingers on the silicon substrate, as well as the interconnection sections between the contact pads on the PW board, and the wires interconnecting the wire bond fingers and the contact pads.\nIn addition, a structural member 730 which acts as a heat sink and encapsulates the cavity section is provided in the device 710. The end parts 731 of the structural member (heat sink) are situated on the upper level section of the PW board. Although the heat sink is located with a spacing from the chips of the MCM tile, it is situated close to the MCM tile to such a degree that it is sufficient to collect heat generated by the constituent elements of the MCM tile during the operation of the device. As an option, a heat conductive adaptive member 732 such as a heat conductive paste or a thermal grease may be given so as to make physical contact with the chips and the heat sink.\nA second example of MCM package 770 is shown in FIG. 7B. This is an example of an MCM tile having a constitution in which it is interconnected to the PW board by solder reflow bonding. The MCM package has a PW board 771 formed of a single level section which has a through opening 772. The positional relation between the PW board 771 and an MCM tile 717 is such that chips 719 and 720 of the MCM tile are within the opening 772, while the end parts of a silicon substrate 718 of the MCM tile are positioned to overlap with the bottom face of the PW board 771 adjacent to the opening so as to have the silicon substrate 718 of the MCM tile to be on the outside of the opening.\nEach of bond fingers 773 on the silicon substrate is connected electrically to a contact 774 on the PW board by solder reflow interconnection. A cup-shaped cover 775 makes contact with the bottom face of the silicon substrate 718 while the flanges 776 of the cover are attached to the bottom of the PW board 771 by means of an adhesive (not shown). In order to use the cup-shaped cover as a heat sink for the MCM tile, it is formed of a metal such as copper or a plastic having a high heat conductive property. In the case of a metallic cover, it has an advantage that it acts as a shielding body against electromagnetic radiation.\nA cavity section 777 is formed by the wall sections of the opening 772 and the cup-shaped cover, an encapsulation sealing material with adaptability such as silicone gel is filled partially, and the sealing material 729 seals and protects the interconnection part between the MCM tile and bond fingers, and the contacts.\nThe conventional MCM semiconductor device has a configuration in which a plurality of IC chips are mounted on a silicon substrate being an intermediate substrate, and the silicon substrate is mounted on a PW board, as described, for example, in the semiconductor device disclosed in Japanese Patent Applications Laid Open, No. Hei 8-250652. Accordingly, the size of the silicon substrate becomes extremely large compared with the size of the IC chips, but no consideration on their size is given there. However, in the configuration in which a silicon substrate is adhered to the entire surface of the PW board, as in the example shown in FIG. 7A, for example, there occurs a problem that the silicon substrate tends to have cracks due to thermal stress when the size of the silicon substrate is increased. Moreover, wirings mutually connecting the IC chips are formed on the silicon substrate along with the wirings for external connection. For a size of the silicon substrate which exceeds, for example, 20 mm×20 mm, there arises a limit at present in the refinement of a connection wiring pattern because it is impossible to form a wiring pattern in one time of exposure treatment, bringing about also a problem that restricts the realization of higher density for the connection wirings.\nNow, by forming an opening in a PW board for mounting an intermediate substrate to accommodate IC chips mounted on the intermediate substrate in the opening, and connects the intermediate substrate to the PW board using only the electrode section provided in the periphery of the intermediate substrate, as in the example in FIG. 7B and the semiconductor device disclosed in Japanese Patent Applications Laid Open, No. Hei 9-181256, the problem of cracks in the intermediate substrate due to thermal stress can be relaxed. However, since a large opening is provided in the central part of the PW board, there still arises another problem that the number of external connection electrodes of the semiconductor device is limited or that it is necessary to enlarge the size of the PW board in order to secure a prescribed number of electrodes."}
{"text": "Phase change memory (PCM) is an emerging segment of semiconductor memory technology. Phase change memory operation requires a variety of DC power supply voltages to support read and write operations, with some voltages being higher than the external power supply. For example, phase change memory operation requires a variety of voltages to be generated on die to support operation with good DC characteristics of up to μS duration, e.g., ˜0.4 V for bitline precharge, ˜1 V for standard logic and ˜2.5 V for wordline, read sense, and write supply. Tight power supply tolerances are required for resistance sensing and writing PCM arrays with differing load currents and duration.\nVoltage charge pump circuits are required to raise the voltages higher than the external power supply. The requirements/demands on the high voltage supply are challenging and include the requirement of low ripple while being able to supply current over a wide range, high current supply capability, and good tolerance. However, voltage charge pump systems developed from other types of memory devices do not serve the unique requirements of phase change memory sufficiently. In other words, they do not supply high voltage over a wide range, with high current capability while also having low ripple.\nCharge pumps are used in integrated circuits to provide a boosted supply voltage in applications such as eDRAM memory, FLASH memory, bandgap voltage references, etc. Typical boost circuits first charge a capacitor from an external supply and then transfer the stored charge into a capacitor on a boosted supply net. A ripple voltage, though, develops on the boosted supply net that needs to be minimized. Currently, these ripples are minimized by placing large filter capacitors (also called decoupling capacitance) on the boosted supply net and by the use of multiphase pumps.\nFor a given design, a charge pump supplies current to drive an intended load. The output load current can vary greatly depending upon the different operating modes or other conditions. Further, the pump output current will increase with higher pump frequency and supply voltage. To regulate the output voltage of a charge pump with good precision and a tight voltage tolerance it is necessary to use a regulator circuit to control a charge pump, typically turning it on and off so at to keep the output voltage within a desired range. Typically, a charge pump provides some excess charge before it can be turned off. In addition, the boosted voltage must fall to a lower potential before a regulator will turn the pump back on. The excess charge and fall in potential is a characteristic of regulation and adds to voltage ripple on the boosted supply net.\nMore specifically, a charge pump system can be sized to provide adequate current at low supply voltage. However, it will also typically have excessive ripple voltage unless the dcap is sized for the higher charge transfer at the maximum supply voltage. But, decoupling capacitance typically consumes close to 50% of the total pump area in the design and has a significant impact on chip size."}
{"text": "The invention relates to hydraulic gear machines, such as pumps or motors, especially internal gear machines or spur gear machines of the type having a first gear which is rotatably journalled in a housing and a second gear which drives the first gear, as well as with at least one pressure building slot which is provided in the housing.\nHydraulic gear machines of the above outlined character are known, for example, from published German patent application Serial No. 2942417, U.S. Pat. No. 4,472,123 and published German patent application Serial No. 1553014. These gear machines have pressure building slots which are provided in the axial plates of the housing and, in a plan view, have a triangular outline. The pressure building slots in the axial plates of the housing in a gear machine serve to ensure a more uniform pressure application or pressure rise which is effected all the way to a plenum chamber and, in the event that the gear machine is used as a pump, there takes place a pressure rise and a rise in the rate of delivery from a suction chamber to a plenum chamber. When such gear machine is in use, it conveys a fluid from the suction chamber to the plenum chamber.\nWhen such gear machines are used as pumps for motor vehicle aggregates, and especially for the transmissions of motor vehicles, such as for continuously variable speed transmissions with pulleys employing conical flanges, for example to convey a fluid, such as oil, from a sump to act upon adjusting members and/or for lubrication and/or cooling, the pump draws from the sump--during certain stages of operation of the vehicle--oil which contains a relatively high percentage of air, for examples a mixture of air and oil whereas, during other stages of operation of the vehicle, there is conveyed oil or an oil-air mixture which contains a relatively low percentage of air.\nIn known prior art pumps, the above brings about the drawback that the preliminary compression or the regulation of pressure by way of pressure building slots which--when looked at from above--have a triangular outline, and when the oil contains a relatively high percentage of air, takes place abruptly, i.e., in response to a relatively small angular displacement of the pinion. On the other hand, when the oil contains a relatively small percenage of air, the regulation of pressure takes place in response to a relatively larger angular displacement of the pinion, namely an angular displacement corresponding to the extension of the pressure building slot.\nPressure equalization during a relatively small angular displacement when the oil contains a high percentage of air entails a rather unsatisfactory acoustic behavior of the pump because, when the percentage of air in oil is relatively high, this causes the pressure rise to take place at a relatively high speed which, in turn, generates pressure waves causing a reduced acoustic quality."}
{"text": "Gas turbine engine components are required to operate efficiently during engine operation and flight. Components within the gas turbine engine, such as gear systems, may aid in the rotation of an engine shaft, turbines and fans at a high speed. Gear systems rotating at a high speed require lubrication in order to operate efficiently and also to reduce damage to the engine. Certain conditions may prevent lubrication to portions of a gas turbine engine.\nAccordingly, there is a need for lubrication systems to maintain adequate lubrication to gear systems of a gas turbine engine."}
{"text": "This invention relates to a solid-state imaging device, and more particularly to a circuit for suppressing leakage current at the signal storage nodes of a CMOS solid-state image sensor (CMOS image sensor) with a read circuit capable of reading an image signal pixel by pixel, a circuit of suppressing the current drawn by an image signal amplifying source follower, and a horizontal readout gate providing a horizontal signal line with less parasitic capacitance, which are used in, for example, a video camera or an electronic still camera.\nFIG. 20 shows an equivalent circuit of a conventional CMOS solid-state image sensor (amplification CMOS image sensor) (conventional equivalent 1) with a read circuit capable of reading an image signal pixel by pixel.\nIn the sensor of FIG. 20, a cell area (imaging area) 1 is composed of plural unit cells 13 arranged in a two-dimensional matrix. One unit cell corresponds to one pixel.\nEach unit cell 13 is composed of, for example, four transistors and one photodiode. Specifically, each unit cell 13 includes a photodiode 8 to whose anode the ground potential is applied, a read transistor (shutter gate transistor) 14 one end of which is connected to the cathode of the photodiode 8, an amplifying transistor 15 whose gate is connected to the other end of the read transistor 14, a vertical select transistor 16 one end of which is connected to one end of the amplifying transistor 15, and a reset transistor 17 one end of which is connected to the gate of the amplifying transistor 15. In the cell area 1, the following lines are formed: read lines connected in common to the gates of the individual read transistors of the unit cells in the same rows, vertical select lines 6 connected in common to the gates of the individual read transistors 14 of the unit cells in the same rows, reset lines 7 connected in common to the gates of the individual reset transistors of the unit cells in the same rows, vertical signal lines 18-i (i=1 to n) connected in common to the other end of the individual amplifying transistors 15 of the unit cells in the same columns, and power lines 9 connected in common to the other end of the individual reset transistors and to the other end of the individual vertical select transistors 16 of the unit cells in the same columns.\nOutside the cell area 1, the following component parts are provided: load transistors 12 connected between one end of the respective vertical signal lines 18-i and the ground nodes, horizontal select transistors 23-i one end of which is connected to the other end of the respective vertical signal lines 18-i via the corresponding noise chancellor circuits 25-i, a horizontal signal line 26 connected in common to the other end of the horizontal select transistors 23-i, an output amplifier circuit 27 connected to the horizontal signal line 26, a horizontal reset transistor 28 connected to the horizontal signal line 26, a vertical shift register 2 for supplying a select signal in a scanning manner to the vertical select lines 6 of each row in the cell area 1 and driving the vertical select transistors 16 in each row in a scanning manner, a horizontal register 3 for driving the horizontal select transistors 23-i in a scanning manner, and a timing generator circuit 10 for generating various timing signals.\nEach of the noise chancellor circuits 25-i is composed of, for example, two transistors and two capacitors. Specifically, each noise chancellor circuit is composed of a sample hold transistor 19 one end of which is connected to the other end of the vertical signal line 18-i, a coupling capacitor 20 one end of which is connected to the other end of the sample hold transistor 19, a charge accumulation capacitor 21 connected between the other end of the coupling capacitor 20 and the ground node, and a potential clamping transistor 22 connected to the junction node of the capacitors 20, 21. One end of the corresponding one of the horizontal select transistors 23-i is connected to the junction node of the capacitors 20, 21.\nEach of the horizontal select transistors 23-i is made up of an NMOS transistor having an active region (SDG region) formed in a p-well selectively formed at the surface of a semiconductor substrate. The p-well is connected to the ground potential.\nFIG. 21 is a timing waveform diagram to help explain the operation of the solid-state image sensor of FIG. 20. Referring to FIG. 21, the operation of the solid-state image sensor of FIG. 20 will be explained.\nThe incident light on each photodiode 8 is converted photoelectrically and the resulting signal charges are accumulated in the photodiodes 8.\nBefore the operation of reading the signal charge, a high reset signal is applied to the reset line 7 for a specific period of time to reset the gate potential of the amplifying transistor 15. The reset transistor 17 is on for the specific period, resetting the gate potential of the amplifying transistor 15 to a desired potential.\nAt the same time, a high select signal is supplied to the vertical select line (address line) 6 selected in a scanning manner by the vertical shift transistor 2. The select signal from the vertical select line 6 turns on the vertical select transistor 16. The power supply line 9 supplies a voltage to the amplifying transistor 15 via the vertical select transistor 16. This causes the source-follower-connected amplifying transistor 15 to output a potential proportional to its gate potential to the corresponding vertical signal line 18-i.\nThere is a variation in the gate potential of the reset amplifying transistor 15. As a result, a variation appears in the reset potential of the vertical signal line 18-i connected to the drain of the amplifying transistor 15.\nTo reset the variation in the reset potential of each vertical signal line 18-i, the sample hold transistor 19 is turned on after the reset transistor 17 has been turned on. As a result, the reset potential of the vertical signal line 18-i is transmitted to the capacitor 21 via the capacitor 20. Thereafter, the potential clamping transistor 22 is kept on for a specific period, fixing the voltage of the junction node of the capacitors 20, 21 at a constant level.\nNext, the read line 4 corresponding to the desired row is selected (or supplied with a high read signal), turning on the read transistor 14. This causes the accumulated charge in the photodiode 8 to be transferred to the gate of the amplifying transistor 15 via the read transistor 14, which changes the gate potential. The amplifying transistor 15 outputs a voltage signal proportional to the amount of change of the gate potential to the corresponding vertical signal line 18-i.\nAs a result, the change in the voltage signal on the vertical signal line 18-i caused by the read operation after resetting has been transferred to the capacitor 21 via the capacitor 20. This removes the noise introduced in the stages before the noise chancellor circuit 25-i, such as variations in the reset potential of each vertical signal line 18-i occurring in the cell area 1.\nAfter the noise removing operations have been carried out, the sample hold transistor 19 is turned off and the vertical select transistor 16 is also turned off. As a result, the unit cell 13 is brought into the unselected state and the cell area 1 is electrically disconnected from each noise chancellor circuit 25-i.\nThen, the horizontal reset transistor 28 is turned on, which resets the horizontal signal line 26. Thereafter, the horizontal select transistors 23-i are turned on sequentially, causing the voltages at the junction nodes (signal storage nodes SN) of the capacitors 20, 21 to be read sequentially. The read-out voltages are amplified by the output amplifier circuit 27, which then outputs the amplified voltages.\nThe above-described noise removing operations are carried out each time one horizontal line is read from.\nIn the prior art, after the series of noise removing operations has been completed, the vertical select line 6 is returned to the low level, turning off the vertical select transistor 16. As a result, the voltage of the vertical signal line 18-i drops to the ground potential via the load transistor 12. At this time, since one end of the vertical signal line 18-i of the sample hold transistor 19 is biased at the same voltage as that of the substrate (e.g., the p-well) in the noise chancellor circuit 25-i, as much leakage current as cannot be neglected develops. The leakage current changes the voltage at the signal storage node SN.\nIn this case, the leakage currents in the individual sample hold transistors 19 differ from one another. This permits the amount of change of the voltage at the signal storage node SN to differ with the leakage current.\nThereafter, the horizontal select transistors 23-i are tuned on sequentially and the signal is read from each horizontal select transistor 23-i. Since the direct-current potentials of the read-out signals differ from each other, image noise, such as vertical lines, takes place due to the variations in the direct-current potential when the output signal of the image sensor is displayed on the screen of the image display unit.\nAdditionally, in the prior art, the potentials of the signal storage nodes SN are clamped in the noise chancellor circuits 25-i, the potentials are clamped at the ground potential by the clamp transistors 22. One end of the clamp transistors 22 and that of the horizontal select transistors 23-i connected to the signal storage nodes SNs are biased at the same voltage as that of the substrate (the p-well in the example). This permits as much leakage current as cannot be neglected to take place at those transistors.\nAs a result, when the horizontal select transistors 23-i are turned on sequentially, the direct-current potentials of the signals of the horizontal select transistors 23-i later selected vary from the direct-current potential of the signal read from the horizontal select transistors 23-i earlier selected (for example, the former get lower gradually). As a result, image noise, such as vertical lines, occurs.\nFIG. 22 shows an equivalent circuit of another conventional amplification CMOS image sensor (conventional equivalent 2).\nIn FIG. 22, in a cell area (imaging area) 1, unit cells 13 are arranged in a two-dimensional matrix. Like the unit cell 13 of FIG. 20, each unit cell 13 in conventional equivalent 2 is composed of a vertical select transistor (row select transistor) Ta, an amplifying transistor Tb, a reset transistor Tc, a read transistor Td, and a photodiode PD. One pixel is made up of one unit cell 13.\nAs in FIG 22, read lines 4, vertical select lines 6, reset lines 7, vertical signal/lines VLIN, and power supply lines 9 are formed in the cell area 1.\nAs in FIG 22, a load transistor TL is connected between one end of the vertical signal lines VLIN and the ground node outside the lower part of the cell area 1.\nLike the noise chancellor circuits 25-i of FIG. 20, noise chancellor circuits 25 are arranged in the horizontal direction outside the upper part of the cell area 1. Each noise chancellor circuit is composed of a sample hold transistor TSH, a potential clamping transistor TCLP, a coupling capacitor Cc, and a charge accumulation capacitor Ct. A horizontal select transistor TH one end of which is connected to the junction node of the capacitor Cc and capacitor Ct is provided for each noise chancellor circuit.\nA horizontal signal line HLIN is connected in common to the other end of each of the horizontal select transistors TH. A horizontal reset transistor (now shown) and an output amplifier circuit (not shown) are connected to the horizontal signal line HLIN.\nFurthermore, outside the cell area 1, the following are provided: a vertical shift register 2 for selecting the vertical select transistors Ta in each row in a scanning manner, a horizontal shift register 3 for driving the horizontal select transistors TH in a scanning manner, a timing generator circuit for generating various timing signals to be supplied to, for example, the noise chancellor circuit 25, a bias generator circuit 11 for supplying a specific bias potential to one end of the potential clamping transistor TCLP of the noise chancellor circuit 25, and a pulse selector 2a for driving the individual rows in the cell area 1 in a scanning manner under the control of the output pulse from the vertical shift register 2.\nIn FIG. 22, the amplifying transistor Tb of each unit cell 13 and the load transistor TL connected to the amplifying transistor Tb via the vertical signal line VLIN constitute a source follower amplifier circuit.\nThe operation of the solid-state image sensor of FIG. 22 is basically the same as that of the solid-state image sensor of FIG. 20 except for operation timing.\nFIG. 23 is a timing waveform diagram to help explain the operation of the solid-state image sensor of FIG. 22.\nIn each unit cell 13, incident light on the photodiode PD is converted photoelectrically. The resulting signal charges are accumulated within the photodiode PD.\nTo read the signal charges in the photodiodes PD from the unit cells 13 in a row (hereinafter, referred to as a selected row) in the horizontal retrace line interval, the signal (φ ADRES pulse) to the vertical select line 6 for the selected row is turned on, turning on the row select transistors Ta in the selected row, to select each of the vertical signal lines VLIN.\nThis causes the source follower circuit composed of the amplifying transistor Tb supplied with the power-supply potential VDD via the row select transistor Ta and the load transistor TL to operate in the individual unit cells 13 in the selected row.\nNext, in each unit cell in the selected row, the signal (φ RESET pulse) to the reset line 7 is turned on to reset the gate voltage of the amplifying transistor Tb to a reference voltage, thereby outputting the reference voltage to the vertical signal line VLIN.\nIn this case, the driving signal (φ SH pulse) of the sample hold transistor TSH in the noise chancellor circuit 25 is turned on in advance (for example, at the same time the φ ADRES pulse is turned on). After the reference voltage has been outputted to the vertical signal line VLIN, the driving signal (φ CLP pulse) of the potential clamping transistor TCLP is kept on for a specific period of time, which sets the reference voltage in the noise chancellor circuit 25.\nNext, after the φ RESET pulse has been turned off, the signal (φ READ pulse) on the read line 4 is turned on. This turns on the read transistor Td, supplying a voltage proportional to the accumulated charges in the photodiode PD to the gate of the amplifying transistor Tb. As a result, the signal voltage proportional to the accumulated charges is outputted to the vertical signal line VLIN and noise chancellor circuit 25.\nThereafter, the φ SH pulse in the noise chancellor circuit 25 is turned off, which allows the signal component (noise-removed signal voltage) corresponding to the difference between the reference voltage and the read-out signal voltage to be accumulated in the charge accumulating capacitor Ct even in the effective horizontal scanning period.\nThen, after the cell area 1 has been electrically separated from each noise chancellor circuit 25, the signal voltage accumulated in the capacitor Ct turns off φ ADRES pulse in the effective scanning period, turning off the vertical select transistor Ta, which places the unit cells 13 in the selected row in the unselected state. Thereafter, sequentially turning on the driving signals (φ H pulse) for the horizontal select transistors TH causes the horizontal select transistors TH to turn on in sequence, allowing the signal voltage accumulated on the horizontal signal line HLIN to be outputted.\nIn the operation, the voltage VVLIN of the vertical signal line VLIN becomes the operating voltage Vm (about 1.5V) of the source follower circuit in the horizontal retrace line interval, whereas it becomes 0V in the effective horizontal scanning period. This permits leakage current from the sample hold transistors TSH to occur in the effective horizontal scanning period and enter the capacitors of the noise chancellor circuit 25. The leakage current differs from one vertical line to another, which causes image noise, such as vertical lines.\nAs described above, the conventional CMOS image sensors have the problem of permitting image noise, such as vertical lines, to appear on the display screen of the image sensor output signal as a result of changes in the voltage at the signal storage node caused by the current leakage after a series of noise removing operations has been completed by the noise chancellor circuit in a read operation for each horizontal line.\nFIG. 25 shows patters of part of the horizontal readout gate section of the conventional CMOS image sensor of FIG. 20.\nIn FIG. 25, reference symbol 23a-i (i=1 to 4) indicates the active region (SDG region) of the horizontal select transistor 23-i formed in a p-well selectively formed at the surface of a semiconductor substrate. Element isolating regions 24 are formed between SDG regions.\nReference symbol 23b-i indicates the gate electrode (polysilicon wire) of the horizontal select transistor 23-i. The gate electrode is formed on the channel of the SDG region 23a-i via a gate insulation film (not shown) formed at the surface of the p-well.\nA metal wire (normally an aluminum wire) corresponding to the vertical signal line 18-i is connected to the n-type diffused region (source) at one end of the SDG region 23a-i. To the n-type diffused region (drain) at the other end, a metal wire (normally an aluminum wire) corresponding to the horizontal signal line 26 is connected.\nSince the coupling capacitance with the p-well is present in the n-type diffused regions acting as the drain and source of the horizontal select transistor 23-i, the parasitic capacitance 29 on the horizontal signal line 26 increases in proportion to the number of the horizontal select transistors 231-i.\nThe increase in the parasitic capacitance on the horizontal signal line 26 makes the operating speed of the circuit slower. Switching noise developing as a result of the switching operation of the horizontal select transistor 23-i is not removed by the noise chancellor circuit 25-i and enters the parasitic capacitance 29. The amount of incoming noise becomes larger as the parasitic capacitance increases. The larger amount of incoming noise contributes the occurrence of image noise, such as vertical lines when the output signal of the solid-state image sensor is displayed on the screen of the image display unit.\nAs described above, the conventional solid-state image sensor has the problem of making the circuit operation speed slower because the parasitic capacitance of the horizontal signal line increases in proportion to the number of horizontal select transistors and of permitting image noise, such as vertical lines, to appear on the display screen of the output signal from the image sensor as a result of noise entering the parasitic capacitance."}
{"text": "1. Field of the Invention\nThe invention pertains generally to fiber optics, and more particularly to a fiber optic connector housing, a fiber optic receptacle, various accessories for electro-optic modules which include fiber optic connector housings and corresponding optical assemblies.\n2. Description of the Related Art\nThose engaged in the manufacture and use of communication systems, e.g., systems for communicating voice, video and/or data, have become increasingly interested in using fiber optic cables as transmission media in such systems. This interest is stimulated by the fact that the potential bandwidth. (or information-carrying capacity) of optical fibers is extremely high. In addition, communication systems employing fiber optic cables are resistant to electromagnetic interference, which sometimes plagues systems employing electrical cables as transmission media. Moreover, communication systems employing fiber optic cables are considered more secure than systems employing electrical cables because it is generally more difficult for unauthorized personnel to tap or access a fiber optic cable without being detected.\nAn exemplary communication system employing a fiber optic cable as a transmission medium is one which includes, for example, two or more computers, e.g., mini-computers, with each adjacent pair of mini-computers being interconnected by a fiber optic cable which includes two optical fibers, i.e., a transmit optical fiber and a receive optical fiber. Obviously, each mini-computer generates and receives information, i.e., data, in electrical form. Consequently, each mini-computer is also provided with an electro-optic module, typically mounted on a printed circuit board or printed circuit card of the mini-computer, which converts the electrical signals generated by the mini-computer into optical signals, which are transmitted to the adjacent mini-computer via the transmit optical fiber. In addition, the electro-optic module converts optical signals communicated to the mini-computer via the receive optical fiber into corresponding electrical signals.\nAn electro-optic module 10, typical of the type referred to above, is depicted in FIG. 1 and includes a housing 20 containing a transmitter optical subassembly (TOSA) 30 (not shown), a receiver optical subassembly (ROSA) 40 (not shown) and a pinned ceramic substrate 50 (not shown) bearing a number of semiconductor integrated circuit devices, with the pins 55 of the ceramic substrate protruding from the housing 10. The TOSA 30, which is electrically connected to certain of the semiconductor integrated circuit devices (hereinafter denoted the TOSA ICs), includes an electro-optic transducer, such as a semiconductor laser, which serves to convert electrical signals generated by the TOSA ICs into corresponding optical signals. The TOSA 30 also includes a lens (not shown) and a hollow cylinder (for the sake of convenience, hereinafter termed a bore) 35, which bore protrudes from the housing 10, with the lens serving to focus the light produced by the semiconductor laser onto the end of a transmit optical fiber which is to be inserted into the bore 35. Similarly, the ROSA 40, which is also connected to certain of the semiconductor integrated circuit devices (hereinafter denoted the ROSA ICs), includes an electro-optic transducer, such as a PIN photodiode, which serves to convert optical signals received by the photodiode into corresponding electrical signals, which are communicated to the ROSA ICs. The ROSA 40 also includes a bore (hollow cylinder) 45, which protrudes from the housing 10, into which a receive optical fiber is to be inserted, the receive optical fiber serving to communicate optical signals to the PIN photodiode.\nThe transmit and receive optical fibers, referred to above, are depicted in FIG. 1 and are denoted, respectively, by the numbers 60 and 70. As is conventional, and as shown in FIG. 1, each of these optical fibers is encased in one or more protective plastic sheaths, and each of these optical fibers extends into a ferrule (not shown). To prevent optical losses, it is important that the ferrule containing the transmit optical fiber 60 be inserted into the bore 35 so as to bring the transmit optical fiber into precise alignment with the semiconductor laser and corresponding lens of the TOSA 30. Similarly, it is important that the ferrule containing the receive optical fiber 70 be inserted into the bore 45 so as to bring the receive optical fiber into precise alignment with the PIN photodiode of the ROSA 40. If, for example, the transmit and receive optical fibers are single mode fibers, then the accuracy of each of these alignments must typically be to within one micrometer or less.\nOne set of devices which permits the achievement of micrometer-accurate alignment is also depicted in FIG. 1. That is, as shown in FIG. 1, such alignment accuracy is achievable by inserting the ferrule containing the transmit optical fiber 60 into a plug frame (not shown) which, in turn, is inserted into an individual fiber optic connector housing (FOCH) 80. Similarly, the ferrule containing the receive optical fiber 70 is inserted into a plug frame (not shown) which, in turn, is also inserted into an individual FOCH 90. As also shown in FIG. 1, each of the individual FOCHs is, for example, of the so-called push-pull type available from NTT (Nippon Telegraph and Telephone Corporation, Tokyo, Japan) and referred to as SC-01 straight plug connector. As depicted, each such SC-01 connector is hollow and generally rectangular in cross-section and has length, width and height dimensions of, respectively, 1.000 inches (25.4 mm), 0.356 inches (9.05 mm) and 0.293 inches (7.45 mm). Each such SC-01 connector is also to be inserted into a common receptacle housing (described below), which serves as the mechanism for achieving the above-described alignment. To achieve proper orientation of the individual FOCHs relative to this common receptacle housing, the individual FOCH 80 includes a key 82 on a side surface, and the individual FOCH 90 includes a key 92 on a side surface, which keys are to be received in corresponding keyways in the common receptacle housing. In addition, for reasons explained below, the front end of the individual FOCH 80 includes symmetrical, inclined surfaces 84 and 85, and the top and bottom surfaces of the FOCH 80 include symmetrical apertures 86 (shown) and 87 (not shown). Similarly, the front end of the individual FOCH 90 includes symmetrical, inclined surfaces 94 and 95, and the top and bottom surfaces of the individual FOCH 90 include symmetrical apertures 96 (shown) and 97 (not shown). As explained below, it has been believed that the symmetries associated with these inclined surfaces and apertures are essential to achieving alignment accuracies of one micrometer or less.\nTo take into account manufacturing tolerances associated with the common receptacle housing and/or the individual FOCHs 80 and 90, while still assuring successful insertion of the individual FOCHs into the common receptacle housing, the devices depicted in FIG. 1 also include a separate adapter 100 of the type disclosed in U.S. Pat. No. 4,953,929, which is hereby incorporated by reference. This adapter 100, when connected to the individual FOCHs 80 and 90, permits successful insertion to be achieved while taking account of manufacturing tolerances because the adapter serves to maintain the individual FOCHs in a substantially side-by-side relationship while permitting the individual FOCHs to move relative to one another in at least four different directions. That is, as shown in FIG. 1, the adapter 100 includes a generally C-shaped clamp member 110, which is adapted to clamp onto individual FOCH 80, and a generally C-shaped clamp member 120, which is adapted to clamp onto individual FOCH 90. Each such clamp member includes tabs 105 intended to engage corresponding slots or openings in the individual FOCHs. In addition, the adapter 100 also includes a generally S-shaped flexible member 130 which extends between the clamp members 110 and 120. It is the clamp members 110 and 120 which serve to maintain the individual FOCHs in a substantially side-by-side relationship. On the other hand, it is the generally S-shaped flexible member 130 which permits the individual FOCHs to move relative to one another in at least four different directions.\nThe relative motions permitted by the generally S-shaped flexible member 130 are depicted in FIGS. 2a through 2h, with the arrows in these figures indicating the directions of the motions. For example, as depicted in FIGS. 2a and 2b, the flexible member 130 permits the individual FOCHs to be moved compressively and expansively toward and away from each other. In addition, as depicted in FIGS. 2c and 2d, the flexible member 130 permits each individual FOCH to be moved up or down relative to the other FOCH. Further, as depicted in FIGS. 2e and 2f, the flexible member 130 permits the individual FOCHs to be pivoted, relative to each other, about parallel axes 140 and 150, which are perpendicular to the plane of the paper containing FIGS. 2e and 2f. Moreover, as depicted in FIGS. 2g and 2h, the flexible member 130 permits each individual FOCH to be pivoted, relative to the other FOCH, about an axis 160 which extends between the individual FOCHs and is perpendicular to the axes 140 and 150.\nIt should be noted that all parts of the adapter 100 are conventionally of identical thickness which is, for example, 0.030 inches (0.76 mm). In addition, the width of each of the clamp members 110 and 120 is conventionally identical and is, for example, 0.276 inches (7.00 mm), while the height of each of the clamp members 110 and 120 is conventionally identical and is, for example, 0.356 inches (9.05 mm). Further, the width of the flexible member 130 is, for example, 0.142 inches (3.60 mm). Consequently, when the clamp members 110 and 120 are clamped onto the individual FOCHs 80 and 90, and if one takes into account the thicknesses of the vertical side walls of the clamp members and the width of the generally S-shaped flexible member 130, then it follows that the center-to-center spacing between the FOCHs, and therefore the center-to-center spacing between the ferrules contained in the FOCHs, is 0.5 inches (12.7 mm).\nThe common receptacle housing, referred to above, is depicted in FIG. 1 and denoted by the number 140. This receptacle housing is generally rectangular in outline and includes two longitudinally extending apertures 150 and 160, which are generally rectangular in cross-section and dimensioned to receive the individual FOCHs 80 and 90. In addition, the aperture 150 includes a keyway 155 adapted to receive the key 82 of the individual FOCH 80, while the aperture 160 includes a keyway 165 adapted to receive the key 92 of the individual FOCH 90.\nAs also shown in FIG. 1, a first pair of clips 170, interconnected by a wall 177, is provided for insertion into the-rear of aperture 150, while a second pair of clips 180, interconnected by a wall 187, is provided for insertion into the rear of aperture 160. The longitudinal extent of the first and second pairs of clips 170 and 180 in the apertures 150 and 160 is chosen in relation to the lengths of the keyways 155 and 165 so that the keyways extend fully to the front ends of the first and second pairs of clips. In addition, the first pair of clips 170 includes symmetrical, inclined camming surfaces 174 and 175 which are adapted to engage the symmetrical, inclined surfaces 84 and 85, as well as the symmetrical apertures 86 and 87, of FOCH 80. Similarly, the second pair of clips 180 includes symmetrical, inclined camming surfaces 184 and 185 which are adapted to engage the symmetrical, inclined surfaces 94 and 95, as well as the symmetrical apertures 96 and 97, of FOCH 90.\nAs is evident from FIG. 1, the wall 177, interconnecting the first pair of clips 170, includes an aperture 178, centrally located between the clips 170, which is adapted to receive bore 35. Similarly, the wall 187, interconnecting the second pair of clips 180, includes an aperture 188, centrally located between the clips 180, which is adapted to receive bore 45.\nWhen assembled, the above-described elements permit the ferrules containing the optical fibers 60 and 70 to extend into the bores 35 and 45 and to be properly aligned to within micrometer tolerances. That is, in the course of assembly, the adapter 100 is used to interconnect the individual FOCHs 80 and 90 by applying the clamp member 110 to the FOCH 80 and the clamp member 120 to the FOCH 90. In addition, the first and second pairs of clips 170 and 180 are inserted into the rear ends of the longitudinally extending apertures 150 and 160 of the receptacle housing 140, and the bores 35 and 45 are then inserted into the centrally located apertures 178 and 188 of the first and second pairs of clips. The front ends of the interconnected FOCHs 80 and 90 are then inserted into the front ends of the longitudinally extending apertures 150 and 160, with the key 82 entering the keyway 155 and the key 92 entering the keyway 165. These keyways 155 and 165, which extend all the way to the first and second pairs of clips 170 and 180, serve to guide the FOCHs 80 and 90 to these pairs of clips and have been believed to be essential to achieving micrometer-accurate alignment.\nWhen the front ends of the FOCHs 80 and 90 contact the first and second pairs of clips 170 and 180, the symmetrically inclined camming surfaces 174 and 175 of the first pair of clips 170 engage the symmetrically inclined surfaces 84 and 85 of the FOCH 80, and the symmetrically inclined camming surfaces 184 and 185 of the second pair of clips 180 engage the symmetrically inclined surfaces 94 and 95 of the FOCH 90. As the FOCHs 80 and 90 are inserted more deeply into the apertures 150 and 160, the camming action effected by the camming surfaces serves to align the FOCHs 80 and 90, and therefore the ferrules contained in these FOCHs, relative to the bores 35 and 45 even as the ferrules enter these bores. As insertion continues, the clips 170 come to extend into the symmetrical apertures 86 and 87 of the FOCH 80 to thereby engage, and grip, the underlying plug frame, while the clips 180 come to extend into the symmetrical apertures 96 and 97 of the FOCH 90 to thereby also engage, and grip, the underlying plug frame. As a consequence, these plug frames and corresponding transmit and receive optical fibers are maintained in proper alignment relative to the bores 35 and 45 and to the corresponding semiconductor laser/lens combination and PIN photodiode associated with, respectively, the TOSA 30 and ROSA 40.\nIt must be emphasized that it has been believed that the above-described assembly procedure achieves micrometer-accurate alignment because the keyways 155 and 165, which extend all the way to the clips 170 and 180, serve to maintain the FOCHs in proper alignment relative to the clips up until the moment the clips engage the FOCHs, and that only such keyways are capable of achieving such alignment. In addition, it has also been believed that micrometer-accurate alignment is achieved because the forces exerted by the clips 170 and 180 on the FOCHs during the insertion procedure are symmetrical, and that such symmetrical forces can only be achieved by employing FOCHs having symmetrically inclined surfaces and apertures. Moreover, it has been believed that any deviation from such symmetrically inclined surfaces and apertures will necessarily lead to asymmetrical forces being exerted on the FOCHs, leading to unacceptably large misalignments between the transmit and receive optical fibers and, respectively, the semiconductor laser/lens combination and PIN photodiode.\nIt should be noted that the assembly procedure, described above, is relatively complex and adds substantially to the cost of the resulting optical assembly, and that there has long been a desire to reduce this cost.\nSignificantly, a number of new communication systems employing fiber optic links have been proposed. One such system, depicted in FIG. 3, serves to connect a plurality of devices in a so-called star network, using fiber optic links. As shown in FIG. 3, these devices include, for example, a direct access storage device (DASD), which includes one or more hard disks. These devices also include, for example, a printer and central processing units (CPUs) contained in, for example, personal computers or engineering work stations. In addition, these devices include, for example, a number of end stations, such as computer terminals, and a terminal concentrator, which allows the computer terminals to communicate with the CPUs in the star network.\nYet another proposed communication system employing fiber optic links is depicted in FIG. 4. In this system, a number of devices are connected to each other, via fiber optic links, through a central optical switch. This system includes, for example, CPUs, a DASD, computer terminals, and a shared memory, consisting of semiconductor memory. In addition, and as depicted in FIG. 4, this system could also include one or more telephone systems.\nA number of standards have been proposed in connection with the communication systems depicted in FIGS. 3 and 4 and described above. For example, the personal computers and/or engineering work stations containing the CPUs mentioned above conventionally have brackets which hold printed circuit cards. Each such bracket also includes an opening or slot which permits the insertion of an input/output cable, such as an electrical cable or a fiber optic cable. If such a personal computer or engineering work station were to be included in one of the communication systems, described above, then this personal computer or engineering work station would have to include an electro-optic module and, for example, a common receptacle housing, of the type described above, mounted on a printed circuit card of the personal computer or engineering work station. A fiber optic connector housing which included, for example, two individual FOCHs 80 and 90 connected by an adapter 100, which individual FOCHs contain transmit and receive optical fibers, would then have to fit into the relevant slot in order to insert the fiber optic connector housing into the common receptacle housing, and thereby connect the corresponding fiber optic cable to the electro-optic module. However, to achieve consistency with other existing standards, a new standard has been proposed by the X3T9.3 committee of the American National Standards Institute (ANSI), which may also be adopted by the International Electrotechnical Commission (IEC), requiring the slot to be 0.293 inches (7.45 mm) high and 0.856 inches (21.75 mm) wide. Moreover, this new proposed standard also requires the center-to-center spacing between the ferrules, and therefore the center-to-center spacing between the individual FOCHs, to be 0.5 inches (12.7 mm). However, while an adapter 100 in combination with two individual FOCHs 80 and 90 can achieve a center-to-center spacing of 0.5 inches, the height of, for example, NTT SC-01 FOCHs is conventionally 0.356 inches (9.05 mm), which is greater than the the height of the slot specified in the proposed standard. As a consequence, such a combination could not be inserted into such a slot. If, on the other hand, the individual FOCHs were to be rotated by ninety degrees, so that the height and width of the individual FOCHs were to be interchanged, and if the adapter 100 were to be manufactured so that the height and width of the adapter were to be interchanged without reducing the width of the flexible member, then a combination of such an adapter and two such individual FOCHs would fit into the slot defined by the proposed standard. Unfortunately, by virtue of the thicknesses of the side walls of the clamp members 110 and 120 of such an altered adapter 100, the center-to-center spacing between the individual FOCHs would be 0,563 inches (14.3 mm), which violates the center-to-center spacing of 0.5 inches required by the proposed standard. It should also be noted that reducing the width of the flexible member, to achieve a center-to-center spacing of 0.5 inches, is not a viable option because any such width reduction significantly degrades the flexibility, and therefore functionality, of the flexible member.\nThus, those engaged in the development of fiber optics have sought, thus far without success, combinations of individual FOCHs and corresponding adapters, and combinations of common receptacle housings and pairs of clips, which require less assembly, and therefore lead to corresponding optical assemblies which are less costly. In addition, those engaged in the development of fiber optics have sought, thus far without success, combinations of at least two individual FOCHs and corresponding adapters which meet the height, width and center-to-center spacing requirements specified in the proposed ANSI standard, discussed above."}
{"text": "Various methods for determining baseband frequency response of an FM transmitter are known. One technique relies upon the FM transmitter response alignment as performed in conjunction with a deviation meter or FM demodulator. These latter devices lend a frequency dependent response characteristic which cannot be quantified independently of the FM transmitter. Another previous technique employs a pre-emphasis network at the input to the FM transmitter having a 6 dB per octave amplitude vs. frequency characteristic curve for altering the modulating signal input to the FM transmitter to compensate a 6 dB per octave decrease in the first sideband output of the FM transmitter as input frequency is increased.\nWhile these prior approaches have proven useful for their intended purposes, they are subject to various objections. Both the FM demodulator and the pre-emphasis methods introduce a frequency dependent response characteristic relative to the particular application. This injected, relative response characteristic is not amenable to an absolute, universal standard. For example, using the pre-emphasis method, FM transmitters at distant locations could be checked against the same quantized standard only by checking the first transmitter with a given pre-emphasis network to the second station for testing of its FM transmitter, or by building exactly the same network for testing at the second station. The first alternative is not practicable, and the second alternative is subject to a certain unavoidable margin of inaccuracy due to manufacturing tolerances, etc. Furthermore, in either alternative, the network is volatile and will drift in its response characteristics.\nThe response of an FM transmitter may be compensated by tuning the response of the corresponding FM receiver or demodulator at the receive end of a transmission system. From a manufacturing and engineering standpoint, there is a need to alleviate this burdensome compensation and instead provide uniformity and interchangability among FM transmitters at the same or different locations, or at different times of manufacture. There is a further need for industry wide compatability of modulated transmission systems. Interchangability is particularly important in standby switching type systems and for customer support in exact replacement parts."}
{"text": "The present invention relates to a tire intended to be mounted on a rim having at least a first frustoconical seat, the generatrix of which has an axially outer end closer to the axis of rotation than the axially inner end.\nSuch a tire is described in international application WO 94/13498. It comprises at least a first bead, intended to be mounted on a first seat inclined towards the outside, the said first bead, having a seat the generatrix of which has its axially outer end closer to the axis of rotation than its axially inner end, and the said generatrix being extended axially towards to the outside by an outer face delimiting the bead toe, the said face forming with the axis of rotation an angle .gamma., open radially and axially towards the outside and of less than 90.degree.. The radial carcass reinforcement of the said tire, anchored in each bead to at least an inextensible annular bead ring, has a meridian profile, when the tire is mounted on its service rim and inflated to its service pressure, with a constant direction of curvature in the sidewall and bead terminated by the toe and which is such that, in said bead, the tangent to the point of tangency of said profile with the inextensible bead ring forms with the axis of rotation an angle .o slashed. open towards the outside of at least 70.degree..\nThe bead of such a tire and more particularly the structure of the anchoring of carcass reinforcement may be variable. In the application WO 95/23073, the radial carcass reinforcement is anchored to the inextensible bead wire by winding around the said wire passing from the heel to the toe of the bead to form an upturn extending in a rubber section, made of rubber mix in the form of a wedge defined by two sides coming from an apex A located beneath the meridian section of the coated bead wire, the radially outer side forming, with a line parallel to the axis of rotation passing through said apex A, an acute angle .o slashed..sub.1, open radially towards the outside and between 20.degree. and 70.degree., and the radially inner side forming with said parallel line an acute angle .o slashed..sub.2, open radially towards the inside, the rubber mix forming the rubber section axially adjacent to the bead wire having a Shore A hardness greater than the Shore A hardness(es) of the rubber mixes located axially and radially above the wire and the rubber section.\nThe combination of the carcass reinforcement meridian profile such as described in the first cited application with a structure of hooking such as described in the second cited application makes possible to obtain a very good compromise between the handling properties of the tire inflated to its recommended pressure and the said properties when the tire is inflated to a reduced pressure or even zero. Whether the initial clamping on the rim be zero or of a given value, the axial displacement of the anchoring bead wire of the carcass reinforcement, which may result from the tension of said reinforcement, which tension is due, for instance, to the inflation pressure of the tire, makes possible, in combination with the axial disposition of the rubber section reinforced by the upturn, to obtain clamping of the toe of the bead on the rim seat. The value of this clamping increases as a function of the tension force to which the carcass reinforcement is subjected, and may become great in the case of high tension, for instance when one of the side walls of the tire is placed under great tension by a transverse force.\nThe preferential solution described in the application WO 95/23073 is such that the carcass reinforcement upturn has a length such that it is in contact with the total perimeter of the rubber section or wedge; it thus forms the two radially outer and inner sides of the rubber section and the side opposite the apex or center of said rubber section, and its end is located axially beyond the point of intersection of the two outer and inner sides. The part of the upturn immediately adjacent to the part of said upturn which is wound around the bead wire may, in the first instance, form the radially outer edge of the rubber section or wedge and then finally the radially inner side of said rubber section, ending beyond the junction point of the two outer and inner sides. It may also form firstly the radially inner side of the rubber section or wedge, then the side opposite the apex of said wedge, and then finally the radially outer side of said rubber section, ending in the same manner as previously described.\nThe above two structures are complicated and difficult to use industrially, and as a consequence, expensive. Furthermore, in severe conditions of load, the contact pressure exerted by the bead toe upon the axially outer hump of the rim and the great temperature that may be reached in the said bead are such that fissures of rubber may occur in the external protecting layer of the bead, fissures which are propagated along the upturn of the carcass reinforcement and reach the radially inner face of the anchoring bead wire with destruction of the bead.\nThe invention has for object to remedy these drawbacks while keeping the excellent properties of unseating of tire beads."}
{"text": "In related art, to interconnect a connection terminal of a semiconductor device and a connection terminal of a substrate, on which to mount the semiconductor device, the practice has been to interconnect the two connection terminals by anisotropic conductive interconnection with the use of an anisotropic conductive interconnection material. With such anisotropic conductive interconnection, a film-like or paste-like anisotropic conductive connection material, obtained on dispersing fine conductive particles in an insulating adhesive, is sandwiched between the two materials to be connected. The resulting assembly is heated and pressured to bond the two materials and the connection material together as electrical conductivity is maintained.\nRecently, with the tendency towards miniaturization and higher performance of semiconductor devices, as exemplified by a liquid crystal display device, miniaturization of circuits as to be used the anisotropic conductive interconnection is progressing. To cope with this, a strong demand is presented for fine pitch between the circuits, while there is felt a strong concern as to the tendency to shorting which may be produced on setting up the anisotropic conductive interconnection. To meet this general demand for the fine pitch, a conductive particle 101 made up of a conductive particle 102, used as an anisotropic conductive interconnection material, and an insulating resin film 103, coated on the particle 102, as shown in FIG. 1, has come to be used to show two functions that are contrary to each other, namely the electrical conductivity and the insulating properties.\nHowever, as the pitch becomes finer, there is presented a problem that, with a conventional insulating film uniformly coated on a conductive particle, if the film thickness is increased to assure the insulating properties between neighboring circuit components, the electrical conductivity is lowered due to the aforementioned characteristics acting in contrary fashion to each other, namely the insulating properties and the electrical conductivity."}
{"text": "Packaging the 3 billion nucleotides of the human genome into the nucleus of a cell requires tremendous compaction. To accomplish this feat, DNA in our chromosomes is wrapped around spools of proteins called histones to form dense repeating protein/DNA polymers known as chromatin: the defining template for gene regulation. Far from serving as mere packaging modules, chromatin templates form the basis of a newly appreciated and fundamentally important set of gene control mechanisms termed epigenetic regulation. By conferring a wide range of specific chemical modifications to histones and DNA, epigenetic regulators modulate the structure, function, and accessibility of our genome, thereby exerting a tremendous impact on gene expression. Hundreds of epigenetic effectors have recently been identified, many of which are chromatin-binding proteins or chromatin-modifying enzymes. Significantly, an increasing number of these proteins have been associated with a variety of disorders such as neurodegenerative disorders, metabolic diseases, inflammation, and cancer. Thus, highly selective therapeutic agents directed against this emerging class of gene regulatory proteins promise new approaches to the treatment of human diseases.\n2-((4S)-6-(4-chlorophenyl)-1-methyl-4H-benzo[c]isoxazolo[4,5-e]azepin-4-yl)acetamide is an inhibitor of one or more bromodomain-containing proteins and is useful in treating bromodomain-containing protein-mediated disorders, such as, e.g., proliferative disorders, inflammatory diseases, sepsis, autoimmune diseases, and viral infections. The amorphous form of 2-((4S)-6-(4-chlorophenyl)-1-methyl-4H-benzo[c]isoxazolo[4,5-e]azepin-4-yl)acetamide is exemplified in U.S. Patent Publication No. 2012/0157428 as Compound 144, and is incorporated herein by reference. 2-((4S)-6-(4-chlorophenyl)-1-methyl-4H-benzo[c]isoxazolo[4,5-e]azepin-4-yl)acetamide is represented by the following structural formula:\n\nGiven the therapeutic benefits associated with the amorphous form of 2-((4S)-6-(4-chlorophenyl)-1-methyl-4H-benzo[c]isoxazolo[4,5-e]azepin-4-yl)acetamide, development of other forms of this compound, having improved properties, represents an attractive area for producing enhanced formulations for inhibition of one or more bromodomain-containing proteins."}
{"text": "This disclosure relates to turbo diesel engines having exhaust gas recirculation.\nLarge diesel engines are used in locomotive and off-road applications. Diesel engines for these applications are in the range of 3-5 megawatts, for example. It is desirable to recirculate the exhaust gases into the charge air stream to reduce the emissions of the diesel engine.\nThe intake pressure of large diesel engines is typically higher than the exhaust pressure. As a result, the exhaust gases must be pumped into the charge air to recirculate the exhaust gases. Electric motors are typically used as the exhaust gas recirculation pump. Electric motors can be difficult to package and may result in a fuel consumption penalty for the overall efficiency of the diesel engine."}
{"text": "Many websites and applications apply various techniques and methods to discern a human user from automated software prior to granting certain access or privileges to the user. Distinction between a human user and automated software is often desired, such as when a user is signing up for a service, when a user is attempting to log in to a service or account, when a user is participating in an online poll or forum, or when a user generating an online post. Many websites employ CAPTCHAs (Completely Automated Public Turing tests to tell Computers and Humans Apart) to achieve such functionality. A common CAPTCHA includes a depiction of distorted text images and requires transcription of the distorted text image. However, websites, applications, and other services that implement such CAPTCHAs risk inadvertently driving away traffic by legitimate users (i.e. non-spamming humans) while providing only minimal protection against automated software and human spammers. Alternatively, websites, applications, and other services that fail to implement some such security feature risk security breaches.\nTherefore, there is a need in the human verification field to create a new and useful method for generating a human likeness score indicative of the legitimacy of a use. This invention provides such a new and useful method."}
{"text": "The field of the disclosure relates generally to a circuit breaker assembly and, more particularly, to high-speed couplings for use in circuit breaker assemblies and methods of assembling such couplings.\nConventional electrical distribution systems employ switching devices, such as for example, circuit breakers in switchgear applications, are conventionally housed in a metal enclosure or “drawout” unit. The drawout units releasably connect the electrical switching devices to an electrical busbar. It is sometimes necessary, such as for scheduled maintenance or testing purposes, to move, or “rack”, the circuit breaker in and out of the enclosure to make and break electrical connection within an electrical power distribution circuit. Conventional drawout units are configured to rack the circuit breakers between a disconnect position, in which the primary contacts or terminals of the circuit breaker are fully disengaged from mating primary contacts within the enclosure, and a connect position in which the primary contacts of the circuit breaker and enclosure are fully engaged. A racking mechanism is typically employed to ensure steady movement of the circuit breaker as it moves between the disconnected and connected positions.\nCircuit breakers are also typically equipped with a variety of auxiliary electronic devices and related control circuitry, including switches, motors, solenoids and the like, for providing control of various functions of the circuit breaker. In such cases, terminals are conventionally arranged on the outside surface of the circuit breaker assembly to provide power to the auxiliary devices and control circuitry, and arranged for automatic connection and disconnection from corresponding mating terminals or contacts arranged within the enclosure. Such corresponding terminals arranged on the circuit breaker and the enclosures are conventionally referred to as “secondary disconnect” terminals or contacts, to distinguish these terminals from the main or “primary” current-carrying circuit breaker and enclosure terminals.\nWhen the circuit breaker is inserted, or racked in, from the disconnect position to the connect position, the corresponding secondary disconnect terminals on the circuit breaker and enclosure must be engaged for power to provided to the auxiliary devices and control electronics. Conversely, when the circuit breaker is removed, or racked out, to the disconnect position, the secondary disconnect terminals must be disengaged. Additionally, it is common to provide a test position, between the connected and disconnected positions of the circuit breaker, wherein the secondary disconnect terminals are engaged, but the primary terminals are not engaged. In such cases, the auxiliary functions of the circuit breaker can be safely tested while the circuit breaker is disconnected from the primary circuit, or “OFF”.\nAt least some known circuit breakers include secondary disconnects rated for between 400-600 volts, and, as such, require relatively large size conductors to accommodate current and larger creepage and clearance spacing to accommodate voltage. As such, at least some known secondary disconnect constructions are not be optimal for CAT5 and CAT6 high speed connections, such as those used in typical Ethernet and HDMI communications, which require signal level voltages typically under 24 volts. Such high speed communications connectors are generally smaller and have less contacting engagement than connectors used for relatively slower speed communications, such as CAT4 communications cables. As such, at least some known circuit breakers position a high speed communications connector outside the main breaker. However, in order to transmit data through the high speed connector, such circuit breakers convert the data into different data types, which slows the data output and prevents the circuit breaker from transmitting the high speed data in real time.\nFor at least the reasons stated above, a need exists for a circuit breaker having a high speed communications connector embedded in the primary electronics to enable transmission of high speed communications data and arranged for automatic connection and disconnection from corresponding mating terminals or contacts arranged within the enclosure."}
{"text": "Canisters are commonly used to hold nuclear fuel in the nuclear power field. Once these canisters have been used and are spent, the canisters need to be transported to large storage tanks positioned in pits in the ground for safety reasons. The process of lifting and transporting spent nuclear fuel canisters requires careful planning and precision. In particular, it is possible for seismic or other movement events to take place during the transport of the canisters, particularly in regions of the world susceptible to earthquakes. Additionally, it is possible to have single failures within a transport system (e.g., a component breaking) during the transport of canisters. Without a system built to handle these types of events, the canisters could tip, fall, be released, and/or otherwise be damaged or compromised."}
{"text": "The invention is directed to a catheter having a basket-shaped electrode array with an improved mechanism for expanding and retracting the electrode array.\nElectrophysiology catheters are commonly-used for mapping electrical activity in the heart. Various electrode designs are known for different purposes. In particular, catheters having basket-shaped electrode arrays are known and described, for example, in U.S. Pat. No. 5,772,590, the disclosure of which is incorporated herein by reference. Such catheters are typically introduced into a patient through a guiding sheath with the electrode array in a folded position within the sheath so that the electrode array does not damage the patient during introduction. Within the heart, the guiding sheath is removed and the electrode array is permitted to expand to be generally basket-shaped. Some basket catheters include an additional mechanism in the form of a wire or the like connected to an appropriate control hand to assist in the expansion and contraction of the electrode array."}
{"text": "Press-hardened steel alloys are being used for sheet-metal parts incorporated in vehicle body structures that may be assembled together with rivets or welding. One example of a press-hardened steel is boron steel sold under the designation Usibor® 22MnB5. Press-hardened steel can be water cooled or oil cooled to a desired level of hardness from 450 to 520 HV. Press-hardened steel may be annealed to reduce the hardness to 140 HV.\nPress-hardened steel parts may be assembled to other steel parts by welding. But, new automotive assemblies may include combinations of parts made of different materials such as aluminum and composite parts. An Ultra High Strength Steel (UHSS) beam formed by press hardening and a composite part or an aluminum part cannot be efficiently joined together in a welding operation. The preferred technique for joining such part assemblies is to rivet or otherwise fasten the parts together. The hardness of such high strength parts poses significant challenges in high volume manufacturing operations because the rivets have difficulty penetrating the UHSS beam."}
{"text": "It is increasingly common for people to use client-side devices to communicate via a radio access network (RAN) with other devices, whether those devices are directly connected to the same RAN or to another network (such as another RAN or a transport network, as examples) to which that RAN directly or indirectly provides access. In a typical arrangement, a base station of a RAN may radiate to define a wireless coverage in which mobile stations (such as cell phones, tablet computers, tracking devices, embedded wireless modules, and other wirelessly equipped communication devices) can operate. In turn, each base station is typically coupled with one or more transport networks, such as the public switched telephone network (PSTN) and/or the Internet for instance. With this arrangement, a mobile station operating within a coverage area of a base station can engage in air interface communication with the base station and can thereby communicate via the base station with various remote network entities or with other mobile stations.\nIn practice, communications over the air interface between a base station and a mobile station are structured in accordance with a particular air interface protocol or “radio access technology,” with communications from the base stations to mobile stations defining a “forward link” (or downlink) and communications from the mobile stations to the base station defining a “reverse link” (or uplink). Numerous such protocols are well known in the art, and others may be developed in the future. Examples of existing protocols include CDMA (e.g., 1×RTT, 1×EV-DO), LTE, WiMAX, iDEN, TDMA, AMPS, GSM, GPRS, UMTS, EDGE, microwave, satellite, MMDS, Wi-Fi (e.g., IEEE 802.11), and Bluetooth. Each protocol may define its own procedures for initiation of calls, handoff between coverage areas, and functions related to air interface communication.\nWithin the context of a RAN, the term “base station” is sometimes used to describe simply a device known as a base transceiver station (BTS), which contains the hardware, antennas, and so forth that actually conduct the over-the-air portion of the communication with the mobile station on behalf of the RAN. At times, however, the term “base station” is used to refer to a combination of (i) one or more BTSs and (ii) a device known as a base station controller (BSC) (or radio network controller (RNC)), which controls the BTS(s) and connects it (them) to the rest of the network. Further, a base station may encompass functions of one or more other RAN entities as well.\nIn a typical scenario, a mobile station registers with a RAN via a particular BTS, and then operates in what is known as “idle mode” on a carrier frequency in a coverage area provided by that BTS. If another telephone, computer, or other communication device attempts to contact the mobile station, the RAN will typically send the mobile station a page message via at least that particular BTS. The mobile station may respond by requesting and establishing communication on what is known as an air interface traffic channel (or simply a traffic channel), and proceed to conduct the relevant communication. In other instances, the mobile station may itself originate a communication, in which case the mobile station sends the RAN an access-request message in order to request and establish communication on a traffic channel to conduct the relevant communication. In either scenario, after the RAN has established a radio link and allocated a traffic channel for use by the mobile station, the mobile station is said to be operating in what is known as an “active mode” or “connected mode”.\nIn some instances, once no substantive communication has taken place between the RAN and the mobile station for a certain amount of time, the RAN will typically reclaim the radio resources (including the traffic channel) that it had allocated to that mobile station, in order to make such resources available for other mobile stations in the coverage area. This transition of the mobile station from (i) having radio resources allocated to it to (ii) not having radio resources allocated to it is often referred to as the mobile station transitioning from active to “dormant”. In some cases, a mobile station may be defined as being dormant if the radio link has been released but a data link layer connection still exists. In such cases, a dormant mobile station may transition back to an active mode by acquiring a radio link and beginning to use the existing data link layer connection. In one example, the amount of time allowed by the RAN to elapse between the last use of the traffic channel and the reclaiming of the traffic channel may be determined based on a value that is specified for a system parameter known as a dormancy timer (or active-to-dormant timer)."}
{"text": "1. Field of the Invention (Technical Field)\nThe present invention relates to an apparatus and method for collecting, transmitting, and distributing solar radiation to the interior of structures or to exterior facilities.\n2. Description of Related Art\nCurrently, there is a need for inexpensive, efficient lighting for structures and facilities, using solar energy.\nThroughout the 20th century, use of the sun as a source of energy has evolved considerably. The sun was the primary source of interior lighting for buildings during the day prior to the 20th century. Eventually, however, the cost, convenience, and performance of electric lamps improved and the sun was displaced as the primary method of lighting building interiors. When solar illumination was no longer used, a revolution in the way buildings, particularly commercial buildings, were designed occurred, making them minimally dependent on natural daylight and almost totally dependent on artificial light. As a result, artificial lighting now represents the single largest consumer of electricity in commercial buildings.\nDuring and after the oil embargo of the 1970s, renewed interest in using solar energy emerged with advancements in systems to introduce daylight into interiors, hot water heaters, photovoltaics, and other types of lighting systems that did not use oil. Today, daylighting approaches are designed to overcome earlier shortcomings related to glare, spatial and temporal variability, difficulty of spatial control, and excessive illumination. In doing so, however, a significant portion of the available visible light is wasted by shading, attenuation, and/or diffusing the dominant portion of daylight, i.e. direct sunlight, which represents more than 80% of the light reaching the earth on a typical day. Furthermore, the remaining half of energy resident in the solar spectrum, i.e. infrared radiation between 0.7 μm and 1.8 μm, is not used by typical daylighting approaches. Additionally, typical approaches add to building heat gain, require significant architectural modifications, and are not easily reconfigured.\nPrevious attempts to use sunlight directly for interior lighting via fresnel lens collectors, reflective light-pipes, and fiber-optic bundles have been plagued by significant losses in the collection and distribution system, ineffective use of non-visible solar radiation, and a lack of integration with co-located electric lighting systems required to supplement solar lighting on cloudy days and at night.\nPrevious attempts at illumination within structures using solar energy have used methods that typically collected the solar energy to charge batteries and to power incandescent or fluorescent lighting, which required electric utility power connection from the residence or business. Electrical power wiring needed to be run from the main utility power supply, thus creating a labor intensive installation process. Running, burying, and connecting electric wire cable is time consuming and requires specialized and skilled labor.\nThe traditional trade-off between night-time illumination energy required and daytime solar energy collected has precluded using only solar energy and has forced inventors to also use main utility power with its inherent complexities of installation.\nThe present invention relates in general to solar energy illumination of the interiors of structures or exterior facilities such as roads or stadiums. The present invention relates more particularly to an illumination system that collects solar energy, transmits the infrared portion through a coil and transmits the visible light portion through one or more transmittal lines, including but not limited to fiber optic cables or optical tubing, first to a directional light pulse delivery system that also is able to control the intensity of the delivered light. The present invention further relates to an illumination system that transmits the visible portion of solar radiation to a discrete directional delivery system that includes delivery to either optical tubes or photovoltaic devices. The present invention further relates to an illumination system that comprises an energy storage system that provides constant illumination to any location during both the day and the night.\nThe present invention comprises a system to collect, transmit, direct, use, and store solar energy during daylight hours. The present invention further comprises light-emitting diodes (LEDs) that are automatically activated when solar radiation is not available, such as at night or during a cloudy day.\nNovel features and further scope of applicability of the present invention will be set forth in part in the detailed description to follow, taken in conjunction with the accompanying drawings, and in part will become apparent to those skilled in the art upon examination of the following, or may be learned by practice of the invention."}
{"text": "An object retrieval system providing an expanse of flexible material variably configured by travel of a material support element between a first position and a second position which generates movement in objects that collect on the material surface for retrieval.\nThere are a variety of circumstances under which objects collect within a particular area and must be retrieved. Retrieval of these objects can be time consuming and expensive due to the number of objects, the amount of objects that collect, or the nature of the typography in which the objects collect. As but one non-limiting example, a significant problem for golf courses can be that golf balls collect in water traps. A conventional manner of retrieving the golf balls and other debris that falls into these water traps can be to periodically hire divers. The divers manually collect the submerged golf balls and debris by hand.\nSpecifically, with respect to utilizing divers to retrieve objects from water traps or pond, a significant problem can be that the process is time consuming and expensive. Divers must be specially trained to work underwater which commands greater pay than for grounds keepers. The golf balls or other debris which collects at the bottom of the water trap or pond must be retrieved by hand. Moreover, this type of work can be more hazardous than routine grounds keeping activities. Additionally, the collection of golf balls or other objects from water traps or ponds can be motivated by the value of the golf balls or other objects which can be recovered. However, collection of golf balls manually from the bottom of a water trap or pond can exceed the value of the recovered golf balls or other ojects.\nAnother manner of retrieving golf balls and debris from water traps in golf course is to submerge a net under the surface of the water and periodically lift the net above the surface of the water with an inflatable device as described in U.S. Pat. No. 6,447,205 or with a lift device such as that shown in U.S. Pat. No. 6,161,988.\nA significant problem with either of these methods can be that the net has fixed securement to a portion of the lift or inflatable device. As the portion of the lift or inflatable device travels from a first position to a second position the net can become unequally tensioned across the surface. This unequal tensioning can lead to a failure of the device to move objects and debris toward the secured perimeter for a variety of reasons. In certain cases, the net can be torn because fixed securement of the lift device to the net creates localized tension in the net that cannot otherwise be released. In other cases, the secured perimeter of the net comes undone and the net loses tension sufficient to maintain a surface on which objects or debris can travel to the perimeter. Finally, the travel of the lift may be restricted to prevent damage to the net resulting in insufficient elevation in portions of the net to allow travel of objects toward the secured perimeter.\nThe instant invention provides an object retrieval system which addresses each of the above-described problems."}
{"text": "1. Field of the Invention\nThis invention relates to a novel (2,3,4,5,6-pentafluorobenzoyl) diphenyl ether compound and a novel fluorine-containing aryl ether ketone polymer.\n2. Description of the Related Art\nThe high-performance fluorinated polymers have been receiving considerable attention as interesting advanced materials for applications as films, coatings for optical and microelectronics devices, gas separation membranes and so on (Cassidy, P. E., Aminabbai, T. M. and Farley, J. M., J. Macromol. Sci.-Rev. Macromol. Chem. Phys. C29 (2 & 3), 365-429 (1989)). The incorporation of fluorine atoms into polymer chains leads to polymers with increased solubility, flame resistance, thermal stability and glass transition temperature while also leading to decreased color, crystallinity, dielectric constant and moisture absorption. Owing these advantages, poly(aryl ether ketone)s (PEK) containing hexafluoroisopropylidene units had been prepared and studied for the use in aerospace and electronic applications (Tullos, G. L. and Cassidy, P. E., Macromolecules, 24, 6059 (1991)). Recently, PEK's containing perfluophenylene moieties were synthesized from perfluorobenzophenone (Mercer, F. W., Fone, M. M., Reddy, V. N., and Goodwin, A. A., Polymer, 38 (8), 1989 (1997)).\nThese polymers, however, are still deficient in solubility and flame resistance.\nAn object of this invention, therefore, is to provide a novel (2,3,4,5,6-pentafluorobenzoyl)diphenyl ether compound for the use as a raw material of a novel fluorine-containing aryl ether ketone polymer exhibiting high mechanical strength and toughness and excelling in electrical properties, thermal oxidative stability, and solubility.\nAnother object of this invention is to provide a novel fluorine-containing aryl ether ketone polymer.\nYet another object of this invention is to provide a fluorine-containing aryl ether ketone polymer exhibiting high mechanical strength and toughness and excelling in electrical properties, thermal oxidative stability, and solubility."}
{"text": "Semiconductor lasers are commonly used as light sources in applications involving optical communications. Such lasers can generate light in response to an application of an electrical signal. A common laser structure is called an “edge-emitting laser” that usually includes a flat junction formed between two pieces of semiconductor material, each having been treated with a different type of impurity. When an electrical current is passed through the laser, light emerges from a plane of the junction region, at an edge of the laser. Examples of such a laser include the Distributed Feedback (DFB) laser, the Fabry-Perot laser, and the Distributed Bragg Reflector Laser (DBR).\nA completed semiconductor laser is generally made in four phases, similar to other semiconductor devices. A first phase involves an extraction and purification of raw semiconductor materials, such as sand for silicon. Next, a method called “crystal growth and wafer preparation” is used to form wafers, which are thin disks of semiconductor material. In a third stage, “chip fabrication” , a plurality of lasers are created in and on the wafer surface. After chip fabrication, the lasers on the wafer must be tested. The electrical test (also called a “wafer sort”) may be part of the last step of chip fabrication or an initial step in the packaging process, which involves separating the wafer into individual lasers (or “chips”) and placing them into protective packages\nA large problem with semiconductor lasers is a high cost of manufacturing associated with them. In order to confirm that a laser is properly functioning, an individual laser must first be cut from the wafer, separated, and at least partially packaged before it can be tested. If a laser is found to be defective, it is discarded. The yield of a wafer may be a key indicator of costs, and often times the yield may be low. Costs become much more problematic especially since the overall costs may include packaging costs for what tests ultimately determine to be a defective chip.\nOne response to this problem, has been the creation of the Vertical Cavity Surface-Emitting Laser (VCSEL), which in contrast to the edge-emitting laser, can emit light from the wafer perpendicular to the plane of the junction region. A light output perpendicular to the surface of the wafer allows VCSELs to be tested at the wafer-level stage by the monitoring of emitted light. In contrast, it is difficult to test an edge-emitting laser at the wafer-level stage, because the laser light is not accessible until the lasers have been cut from the wafer, since the individual placement of the individual lasers on the wafer does not provide sufficient room to insert light monitoring devices between each laser.\nIn summary, the horizontal emission of light by edge-emitting lasers make it difficult to perform testing at the wafer-level stage. As a result, edge-emitting lasers must be partially packaged before test, resulting in additional costs for the manufacturer that are eventually passed on to consumers. Competing lasers, such as the VCSEL, have a testing advantage over the edge-emitting laser, but cannot replace the edge-emitting laser in many applications."}
{"text": "The compound 8-fluoro-2-{4-[(methylamino)methyl]phenyl}-1,3,4,5-tetrahydro-6H-azepino[5,4,3-cd]indol-6-one (“Compound 1”)\nis a small molecule inhibitor of poly(ADP-ribose) polymerase (PARP). Compound 1, and methods of making it, are described in U.S. Pat. Nos. 6,495,541; 6,977,298; 7,429,578 and 7,323,562. Certain salts and polymorphs thereof, of Compound 1, are disclosed in U.S. Pat. No. 7,268,126 and in International Patent Publication No. WO 04/087713. Other publications describing Compound 1 and uses thereof include U.S. Patent Application Publication No. 2006-0074073, and U.S. Pat. Nos. 7,351,701 and 7,531,530.\nPARP is a family of nuclear enzymes responsible for ADP-ribosylation (a post-translational protein modification) in which poly(ADP-ribosyl)transferases transfer the ADP-ribose moiety from NAD+ onto specific amino acid side chains on nuclear target proteins such as histones and DNA repair enzymes and/or onto previously attached ADP-ribose units. In humans the PARP family encompasses 17 enzymes of which PARP-1 is the best-characterized (Otto H, Reche P A, Bazan F et al, In silico characterization of the family of PARP-like poly(ADP-ribosyl)transferases (pARTs), BMC Genomics 2005; 6:139). Pharmacology studies have shown that Compound 1 is an inhibitor of PARP-1 (Ki=1.4 nM) and PARP-2 (Ki=0.17 nM).\nPARP-1 is involved in DNA homeostasis through binding to DNA breaks and attracting DNA repair proteins to the site of DNA damage. PARP-1 through the addition of ADP-ribose units on target proteins provides the energetic resources necessary for chromatin relaxation and the DNA repair process. These actions promote and facilitate DNA repair. Depending on the extent of DNA damage PARP-1 activation and subsequent poly(ADP-ribosyl)ation mediate the repair of damaged DNA or induce cell death. When DNA damage is moderate, PARP-1 plays a significant role in the DNA repair process. Conversely, in the event of massive DNA damage, excessive activation of PARP-1 depletes the cellular ATP pool, which ultimately leads to cell mortality by necrosis (Tentori L, Portarena I, Graziani G, Potential applications of poly(ADP-ribose) polymerase (PARP) inibitors, Pharmacol Res 2002; 45:73-85).\nIn cancer therapy, many useful drugs as well as ionizing radiation exert their therapeutic effect through DNA damage. Enzyme-meditated repair of single- or double-strand DNA breaks is a potential mechanism of resistance to radiotherapy or cytotoxic drugs whose mechanism of action depends on DNA damage. Inhibition of DNA repair pathway enzymes is thus a strategy for the potentiation of anticancer agents. Inhibition of PARP-1 has shown to potentiate the activity of DNA-damaging agents and ionizing radiation in vivo and in vitro. Accordingly, PARP has been identified as a therapeutic target for cancer therapy in combination with DNA damaging agents. (Tentori L, Leonetti C, Scarsella M, et al, Systemic administration of GPA 15427, a novel poly(ADP-ribose) polymerase-1 inhibitor, increases the antitumor activity of temozolomide against intracranial melanoma, glioma, lymphoma, Clin Cancer Res 2003; 9:5370-9. Satoh M S, Poirier G G, Lindahl T, NAD(+)-dependent repair of damaged DNA by human cell extracts, J Biol Chem 1993; 268:5480-7.)\nIn addition to the potential role as chemopotentiator or radiosensitizer agents, more recent evidence has emerged of sensitivity of cell lines, homozygous for either the BRCA1 or BRCA2 mutation, to a PARP inhibitor alone. (Bryant H E, Schultz N, Thomas H D, et al, Specific killing of BRCA-2 deficient tumors with inhibitors of poly(ADP-ribose) polymerase, Nature 2005; 434:913-7. Farmer H, McCabe N, Lord C J, et al, Targeting the DNA repair defect in BRCA mutant cells as a therapeutic strategy, Nature 2005; 434:917-21.) Preliminary clinical data from a Phase I study with a single-agent PARP inhibitor has recently been published (Yap T A, Boss D S, Fong M, et al, First in human phase I pharmacokinetic (PK) and pharmacodynamic (PD) study of KU-0059436 (Ku), a small molecule inhibitor of poly ADP-ribose polymerase (PARP) in cancer patients (p) including BRCA 1/2 mutation carriers, (J Clin Oncol 2007; 25 (Supplement June 20):3529).\nIt is desirable to have crystalline salts and polymorphic forms thereof that possess properties amenable to reliable formulation and manufacture."}
{"text": "Wireless networks are becoming more and more popular. With the popularity of wireless networks increasing, users also demand broader coverage from a wireless network, such as voice, video and data communication support. With the broader demands placed on wireless networks, more complex mechanisms are required for differentiated services.\nOne solution for meeting the broader demands of wireless networks resulted in an increased interest in mesh networking in the home. A typical audio/visual cluster in the home could include a collection of home entertainment appliances that need to connect with each other. Home entertainment appliances can include an instance in the living room, an HDTV, stereo speakers, possibly an amplifier, and a set top box for distribution. Another example is the PC cluster; such a cluster can consist of a PC with a DVD Jukebox, a tuner, a second DVD R/W, and a printer.\nIn metropolitan areas, the use of wireless access points is growing. Already the market is seeing the emergence of informal federation of access points to constitute larger coverage mesh networks. Each access point deployed can be considered a candidate node in proximity mesh networks, whether at the metropolitan scale or the scale of a homeowner's residence. These examples of mesh networking require reliable connectivity among the network nodes. Rather than the cumbersome wired connections, wireless connectivity with auto configuration results in a much better user experience and therefore an accelerated growth of a mesh networking footprint.\nWireless connectivity is an essential element to low cost, ease of use mesh networking. A promising wireless technology to enable the connectivity of these clusters is Wi-Fi. Wi-Fi is a technology that operates in the unlicensed part of the radio spectrum and is based on the existing IEEE 802.11 Standard. The cost of Wi-Fi devices is dropping to become a commodity element in consumer electronics devices as well as an integral part of personal computers.\nThere is a critical issue that must be solved for wireless mesh networking to be a viable scenario: how to grow a mesh network which inherently is an unmanaged network in the traditional sense. Mesh networks tend to be “best effort” deployment, therefore locating and connecting to other nodes within a certain proximity requires a special act on the part of the access point which constitutes a node in an existing mesh network.\nA problem with mesh networks is that changes in local connectivity, even brief incidental changes, such as a dog walking in front of an access point, can result in disruptions in connectivity. If only a small number of nodes are present in the network, this disruption in connectivity can interrupt some nodes ability to reach the greater internet."}
{"text": "The invention concerns a process for the production of LCDs (liquid crystal display). Those LCDs can be in particular PDLCDs (polymer dispersed LCDs) which are to be produced for example in a roll-to-roll process.\nThe production of PDLCDs in a roll-to-roll process has hitherto been obstructed by virtue of the fact that the functional layer of the respective PDLCD which comprises a porous polymer, that is to say support, matrix filled with LC material, is not closed at the surface. After hardening of the matrix the LCs are present in a non-crosslinked, fluid state. By virtue of the fluid LCs at the surface of the matrix, it has hitherto not been possible to carry out a usually multi-stage roll-to-roll process because the functional layer would stick to other surfaces which are brought into contact therewith. That would result in damage to the functional layer."}
{"text": "The invention relates to a detection method and kit for the identification and quantification of analytes in biological specimens by non-captive substrate liquid phase immunoassay techniques.\nObservations of colloidal gold or silver concentrations have been used in immunoassays in conjunction with solid phase diffusion assays. For example, European Patent No. 207,152 discloses a solid phase diffusion assay utilizing a porous sheet having ligands or receptors prebound to the sheet prior to the application of an analyte and a colloidal gold or silver labeled ligand or receptor. However, assays performed on ligand or antiligand coated porous films must be specifically tailored for a particular analyte and, in view of the limited surface area and reaction time with antigens, this process may be ineffective if the affinity of the substrate bound antiligand for the labeled analyte is low. Furthermore, specific areas of the porous sheet substrate may become overloaded with analyte due to restricted radial diffusion, thereby allowing colloidal metals already bound to analytes to pass through the porous substrate without binding to the surface of the porous substrate for subsequent visual detection.\nU.S. Pat. No. 4,853,335 to Olsen et al., discloses a sandwich immunoassay method by premixing a biological specimen with: (1) a colloidal gold labeled ligand or antiligand and (2) solid phase captive particles coated with a ligand or antiligand and applying the subsequent mixture onto the surface of a porous film. By requiring solid phase captive particles this process may distort the visual detection measurements because uncoupled captive particles may block the pores of the substrate and prevent rapid passage of uncoupled colloidal gold.\nIn both of the foregoing patents a ligand, antiligand or receptor is necessarily prebound to a captive substrate in the form of an insoluble porous sheet or a solid captive particle such as a latex particle. It is desirable to develop a method for quantitative and qualitative detection of an analyte, such as C-reactive protein, using a non-captive substrate in order to accelerate passage of non-complexed immunoassay constituents through the pores of the test support. The use of a non-captive substrate technique presents the most simplistic and flexible approach to the rapid and efficient identification and quantification of ligands and analytes and diagnosis of a disease.\nBy eliminating the need for prebinding ligands to solid supports or substrates as a part of the assay system, non-captive substrate techniques promote accurate detection of analytes and cost savings in the production and subsequent use of diagnostic test kits.\nThe object of the present invention is to overcome the deficiencies of prior art solid phase captive immunoassay techniques through the development of a non-captive substrate immunoassay that is easily conducted and accurate in detection and measurement.\nThe present invention is an immunoassay performed on a liquid biological specimen, e.g. serum, plasma, whole blood, etc., for a specific analyte by forming a test sample by contacting the liquid biological specimen with: (a) a binding substance of a ligand, antiligand or receptor capable of binding the analyte and (b) a detector substance of a colloidal metal labeled ligand or antiligand to form a test sample. Both the binding substance and the detector substance are in liquid phases. The test sample is then applied onto a defined zone of an insoluble porous support film having a pore size impassable to a complex formed between the analyte, if present, with the binding substance and the detector substance, but passable to the binding substance and detector substance while remaining uncomplexed in the absence of the desired analyte.\nIf the analyte is present in the test specimen, the analyte binds with both the detector substance and the binding substance to form a visually discernable precipitable complex on the surface of the porous support film. After application of the test sample to the porous support, the surface of the porous support is visually inspected for color to determine the presence and approximate quantity or the absence of the analyte being assayed."}
{"text": "Many types of devices have been developed over the years for the purpose of converting liquids or aerosols into gas-phase fluids. Many such devices have been developed, for example, to desalinate water so as to remove excess salt and other minerals from water. Saline water, or salt water, generally contains a significant concentration of dissolved salts. Seawater has a salinity of roughly 35,000 ppm, or 35 g/L. Seawater is not potable nor suitable for irrigating crops.\nWater may be desalinated in order to be converted to fresh water suitable for human consumption or irrigation. Large-scale desalination typically uses large amounts of energy as well as specialized, expensive infrastructure. As such, it is very costly to use desalinated water instead of fresh water from rivers or groundwater.\nThree methods of desalination include vacuum distillation, reverse osmosis and multi-stage flash.\nIn vacuum distillation, water is boiled at less than atmospheric pressure. Boiling of a liquid occurs when the vapor pressure equals the ambient pressure and vapor pressure increases with temperature. Due to the reduction in temperature, energy is saved.\nReverse osmosis technology involves semi-permeable membranes and pressure to separate salts from water. Less energy may be used than thermal distillation. However, desalination remains energy intensive."}
{"text": "1. Field of the Invention\nThe present invention relates to down-conversion of electromagnetic (EM) signals. More particularly, the present invention relates to down-conversion of EM signals to intermediate frequency signals, to direct down-conversion of EM modulated carrier signals to demodulated baseband signals, and to conversion of FM signals to non-FM signals. The present invention also relates to under-sampling and to transferring energy at aliasing rates.\n2. Related Art\nElectromagnetic (EM) information signals (baseband signals) include, but are not limited to, video baseband signals, voice baseband signals, computer baseband signals, etc. Baseband signals include analog baseband signals and digital baseband signals.\nIt is often beneficial to propagate EM signals at higher frequencies. This is generally true regardless of whether the propagation medium is wire, optic fiber, space, air, liquid, etc. To enhance efficiency and practicality, such as improved ability to radiate and added ability for multiple channels of baseband signals, up-conversion to a higher frequency is utilized. Conventional up-conversion processes modulate higher frequency carrier signals with baseband signals. Modulation refers to a variety of techniques for impressing information from the baseband signals onto the higher frequency carrier signals. The resultant signals are referred to herein as modulated carrier signals. For example, the amplitude of an AM carrier signal varies in relation to changes in the baseband signal, the frequency of an FM carrier signal varies in relation to changes in the baseband signal, and the phase of a PM carrier signal varies in relation to changes in the baseband signal.\nIn order to process the information that was in the baseband signal, the information must be extracted, or demodulated, from the modulated carrier signal. However, because conventional signal processing technology is limited in operational speed, conventional signal processing technology cannot easily demodulate a baseband signal from higher frequency modulated carrier signal directly. Instead, higher frequency modulated carrier signals must be down-converted to an intermediate frequency (IF), from where a conventional demodulator can demodulate the baseband signal.\nConventional down-converters include electrical components whose properties are frequency dependent. As a result, conventional down-converters are designed around specific frequencies or frequency ranges and do not work well outside their designed frequency range.\nConventional down-converters generate unwanted image signals and thus must include filters for filtering the unwanted image signals. However, such filters reduce the power level of the modulated carrier signals. As a result, conventional down-converters include power amplifiers, which require external energy sources.\nWhen a received modulated carrier signal is relatively weak, as in, for example, a radio receiver, conventional down-converters include additional power amplifiers, which require additional external energy.\nWhat is needed includes, without limitation:\nan improved method and system for down-converting EM signals;\na method and system for directly down-converting modulated carrier signals to demodulated baseband signals;\na method and system for transferring energy and for augmenting such energy transfer when down-converting EM signals;\na controlled impedance method and system for down-converting an EM signal;\na controlled aperture under-sampling method and system for down-converting an EM signal;\na method and system for down-converting EM signals using a universal down-converter design that can be easily configured for different frequencies;\na method and system for down-converting EM signals using a local oscillator frequency that is substantially lower than the carrier frequency;\na method and system for down-converting EM signals using only one local oscillator;\na method and system for down-converting EM signals that uses fewer filters than conventional down-converters;\na method and system for down-converting EM signals using less power than conventional down-converters;\na method and system for down-converting EM signals that uses less space than conventional down-converters;\na method and system for down-converting EM signals that uses fewer components than conventional down-converters;\na method and system for down-converting EM signals that can be implemented on an integrated circuit (IC); and\na method and system for down-converting EM signals that can also be used as a method and system for up-converting a baseband signal."}
{"text": "This invention relates to a toy vehicle and, more particularly, to a self-propelled toy vehicle in the shape of a robot which is propelled while performing predetermined tasks.\nIn the field of robot-like toy vehicles, there is an endless demand for novel performance. The present invention is directed to satisfying this demand."}
{"text": "The present disclosure generally relates to cartridges for dispensing a fluid and an automatic analyzer for dispensing the fluid using the cartridge.\nIn medical laboratories, in vitro diagnostics are commonly performed on biological samples such as blood, urine, blood plasma and saliva. Such tests may be performed manually using pipettes or maybe performed using an automatic analyzer. Automatic analyzers may automatically add reagents to the biological sample and may measure one or more parameters of the biological sample during analysis."}
{"text": "This invention relates to a consolidated user interface systems and methods and in particular, to an interface for operating multiple medical systems.\nMany of today's most advanced and effective diagnostic and therapeutic procedures involve the coordinated use of multiple separate medical systems. Each of these separate medical systems may include its own visual display and its own input device. As a result, the visual displays and input devices can take up valuable space in the procedure room, as they must be placed in position for convenient access and use. These multiple displays and input device also pose a risk of confusion among the various displays and input devices."}
{"text": "International Publication No. WO2010/146738 discloses a shift register for scanning gate bus lines. This shift register is composed of a plurality of cascade-connected stages the number of which corresponds to the number of the gate bus lines. Each stage has terminals to which a high-level power source voltage and a low-level power source voltage are input, respectively, and terminals to which an output signal sent from the previous stage and an output signal sent from the next stage are input, respectively. Further, each stage has a terminal to which a clock signal is input, and a terminal through which an output signal of the current stage itself is output as a gate pulse. Each stage is composed of a plurality of transistors including a transistor (M5) and a transistor (M6), and a capacitor that are connected to one another. The transistor (M5) is connected to the terminal to which the high-level power source voltage is input. Further, the transistor (M6) is connected to a terminal into which the low-level power source voltage is input.\nAn internal node in a shift register in each stage has a potential that changes with a clock signal. According to a change of the potential of the internal node, the transistor (M5) is turned on, and a power source voltage VDD is output as a gate pulse via the transistor (M5). This shifts a gate bus line into a selected state. Further, the transistor (M6) is turned ON in response to the clock signal, and the transistor (M6) causes the potential of the internal node to have a power source voltage VSS, thereby decreasing the gate pulse. This shifts the gate bus into a non-selected state."}
{"text": "The present invention relates to a clock signal generator using fractional division of frequency and to means for keeping said clock signal generator under the control of a synchronization signal.\nThe majority of electronic circuits are now designed to operate using a determined frequency clock signal produced by an oscillator. When a piece of equipment includes circuits that are designed to operate on different clock signals, a first solution consists in using a specific oscillator circuit for each clock signal. This solution is clearly not an economical way of solving the problem. Use is often made, when the frequency of a first clock signal is not a multiple of the frequency of a second one, of a clock signal generator employing fractional division of its frequency. Such a clock signal generator, when it is fully implemented in digital technology produces a signal the average frequency of which over one full period is the frequency of the signal it is receiving divided by a division ratio.\nThe use is known of a programmable counter receiving a signal at a steady rhythm and adapted to periodically output, depending on the period of recurrency, a successive first and second series of pulses the frequencies of repetition of which are submultiples of the frequency of the steady rhythm signal. Such a device, apart from the programmable counter, requires two additional counters in order to determine the number of pulses in each series, together with control means. It hence occupies, where it is implemented in integrated form, a fairly large substrate surface area which is an obstacle to the miniaturization of the equipment in which it is incorporated.\nEuropean Patent 0 019 412 disclosed the use of a counter that is preceded by a pulse absorber or subtractor for fulfilling the function of a programmable counter. The pulse subtractor in effect eliminates pulses from the rhythm signal when commanded so to do thus modifying the frequency of the output signal from the counter.\nNevertheless, the phase shift between the output signal and a signal the frequency of which is exactly at the division ratio of the rhythm signal frequency is highly variable and can have an appreciable amplitude.\nA first object of the present invention is to provide a clock signal generator using fractional division of frequency of reduced size made possible by the use of individual circuits that are more simple and the output signal of which only suffers from limited phase shifting relative to the desired frequency. This phase shift, in a preferred embodiment, is less than one period of the steady rhythm signal.\nIn certain fields of application, particularly in telecommunications systems receiving equipment, it is necessary to set the output signal from the clock signal generator under the control of a synchronization signal.\nThe present invention hence also has the object of providing control means that are specially suitable for this type of clock signal generator using fractional division of frequency.\nAlthough one should not attempt to see a limitation of the invention in this, the provisions of the invention can advantageously be used in the Pan-European digital cellular radiocommunications system. In this system, a mobile terminal comprises a speech coding and decoding unit which can be implemented using a standard component working with 8 kHz and 2048 kHz clock signals. The terminal further includes an internal clock signal generator producing a 13 MHz steady rhythm signal, and produces a 50 Hz synchronization signal onto which the coding and decoding units require to be synchronized. The invention hence makes it possible, starting from this steady rhythm signal, to produce clock signals that are synchronized onto this synchronization signal."}
{"text": "The present invention relates to a new and distinctive soybean variety designated XB32AK13, which has been the result of years of careful breeding and selection in a comprehensive soybean breeding program. There are numerous steps in the development of any novel, desirable soybean variety. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The breeder's goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include, but are not limited to: higher seed yield, resistance to diseases and/or insects, tolerance to drought and/or heat, altered fatty acid profile(s), abiotic stress tolerance, improvements in compositional traits, and better agronomic characteristics.\nThese product development processes, which lead to the final step of marketing and distribution, can take from six to twelve years from the time the first cross is made until the finished seed is delivered to the farmer for planting. Therefore, development of new varieties is a time-consuming process that requires precise planning, efficient use of resources, and a minimum of changes in direction.\nSoybean (Glycine max) is an important and valuable field crop. Thus, a continuing goal of soybean breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties.\nThe soybean is the world's leading source of vegetable oil and protein meal. The oil extracted from soybeans is used for cooking oil, margarine, and salad dressings. Soybean oil is composed of saturated, monounsaturated, and polyunsaturated fatty acids. It has a typical composition of 11% palmitic, 4% stearic, 25% oleic, 50% linoleic, and 9% linolenic fatty acid content (“Economic Implications of Modified Soybean Traits Summary Report”, Iowa Soybean Promotion Board & American Soybean Association Special Report 92S, May 1990). Changes in fatty acid composition for improved oxidative stability and nutrition are also important traits.\nIndustrial uses for processed soybean oil include ingredients for paints, plastics, fibers, detergents, cosmetics, and lubricants. Soybean oil may be split, inter-esterified, sulfurized, epoxidized, polymerized, ethoxylated, or cleaved. Designing and producing soybean oil derivatives with improved functionality, oliochemistry, is a rapidly growing field. The typical mixture of triglycerides is usually split and separated into pure fatty acids, which are then combined with petroleum-derived alcohols or acids, nitrogen, sulfonates, chlorine, or with fatty alcohols derived from fats and oils.\nSoybean is also used as a food source for both animals and humans. Soybean is widely used as a source of protein for animal feeds for poultry, swine, and cattle. During processing of whole soybeans, the fibrous hull is removed and the oil is extracted. The remaining soybean meal is a combination of carbohydrates and approximately 50% protein.\nFor human consumption soybean meal is made into soybean flour which is processed to protein concentrates used for meat extenders or specialty pet foods. Production of edible protein ingredients from soybean offers a healthy, less expensive replacement for animal protein in meats and dairy products."}
{"text": "The claimed invention relates generally to the field of data handling systems and more particularly, but not by way of limitation, to a method and apparatus for assigning logical track addresses in a disc drive.\nA disc drive is a data handling system used to store digital data. A typical disc drive includes a number of rotatable magnetic recording discs which are axially aligned and mounted to a spindle motor for rotation at a high constant velocity. A corresponding array of read/write heads are supported by a rotary actuator and used to access fixed sized data blocks (sectors) on tracks of the discs to write data to and to read data from the discs.\nDisc drives are provided with servo control circuitry to move the heads to the various tracks, read/write channel circuitry to write data to and read data from the discs, and interface control circuitry to facilitate communication and data transfer with a host device. A disc drive is typically configured to operate in accordance with an industry standard interface protocol, such as Small Computer Systems Interface (SCSI). Communications and data transfers are carried out between host and drive in accordance with the designated protocol.\nThe available data storage of a typical disc drive is identified at the host level in terms of a sequence of consecutively numbered logical block addresses (LBAs). Each LBA corresponds to a unique data block (sector) at a physical location within the disc stack. A present generation drive can have several million LBAs available to store data.\nDuring a data write command in which the host device writes a computer file to the disc drive, the host operating system uses a file allocation table (FAT) to identify a corresponding number of LBAs across which the file is to be distributed. The host transfers the data to be written to a data buffer of the drive and identifies the target LBAs to which the data are to be written. The disc drive determines the physical locations of the target LBAs in terms of physical tracks and sectors, moves the appropriate head or heads to the respective tracks, and proceeds to write the data to the appropriate sectors.\nA data read command is carried out in a similar fashion; once the host determines need for a previously stored computer file, the host checks the FAT to identify the LBAs in which the file is stored and instructs the disc drive to retrieve the data from the associated LBAs. The drive identifies the physical sectors associated with the target LBAs, schedules movement of the head or heads to the appropriate sectors, and transfers the data to the buffer and then on to the host.\nOf particular interest is the manner in which the disc drive identifies the physical locations of the target LBAs during a read or write (data access) operation. To explain this more fully, it will be helpful to first review the manner in which tracks are defined and addressed in a typical disc drive.\nTracks are defined by servo data recorded to the discs during disc drive manufacturing. The servo data are typically arranged in a number of radially spaced apart servo data fields. Data sectors are subsequently defined in regions between adjacent servo fields during a disc drive formatting operation. The servo data fields on each track typically include a track address field that stores a physical track address (PTA) for that track, typically in Gray code format. All of the tracks at a given radius (i.e., all the tracks having the same PTA) make up a cylinder.\nLogical track addresses (LTAs) are assigned to the tracks over the user data recording areas of the disc surfaces. It is common to leave relatively small guard bands of tracks near the innermost and outermost diameters of the disc surfaces to allow storage of control parameters and information as well as to provide operational margin for the drive. An exemplary LTA scheme involves assigning the first logical track address (logical track zero) to all of the tracks in a first cylinder near the outermost diameter of the discs (such as at PTA=50) and then incrementing the logical track addresses across the discs toward the innermost diameter. Thus, logical track one would correspond to physical track 51, and so on.\nWhen an access command is received, the interface circuit consults a conversion table or otherwise performs the necessary calculations to identify head, logical track and sector for each LBA associated with the access command. The interface circuit directs the servo control circuitry to carry out a seek to the logical track.\nIn order to do so, the servo control circuitry converts the logical track address to a physical track address; using the example above, a command to move a selected head to logical track 0 is interpreted by the servo control circuitry to require movement of the head to physical track 50 (in terms of Gray code). The servo circuit proceeds to execute a seek to move the selected head to the target track. Once the servo control circuit reports that the head is on track, the interface circuit directs the read/write channel to read or write the data to the target sector(s) on that track.\nAs will be appreciated, computer files are often much larger than a single sector and a given data access operation can involve accessing multiple sectors on multiple disc surfaces. This can be especially true in high throughput sustained accesses such as video streaming. While the heads are nominally aligned, small radial offsets (variations in radial location) will typically occur from head to head due to a number of factors such as manufacturing variations, deflection during handling, etc. Thus, if a first head is over physical track 50 and a head switch is made to a second head having a xe2x88x922 track offset with respect to the first head, the second head will actually be over physical track 48 when the head switch operation is completed. A two track corrective seek will be necessary to move the second head to cylinder 50 to continue the data access operation.\nIt is generally known in the prior art to measure head offsets and account for such during servo operation. Using the above example, if the first head is over physical track 50 and the servo circuit is directed to move the second head over physical track 60, the second head will actually need to move 12 tracks (not just 10) to get to the desired track 60, and adjustments can be made accordingly. However, during an access operation to a given logical cylinder where multiple physical tracks are sequentially accessed by different heads, knowing the head-to-head offsets does nothing to eliminate the need to proceed with small seeks after head switch operations to maintain the heads over the tracks in the selected logical cylinder.\nAs track densities continue to increase, the magnitudes of corrective seeks (in terms of tracks) after head switch operations will continue to increase and such seeks will generally tend to decrease effective data transfer rate performance of a disc drive. Accordingly, there is a need for improvements to address these and other limitations of the prior art.\nIn accordance with preferred embodiments, a data handling system (disc drive) is provided in communication with a host device. The disc drive has a plurality of transducing heads adjacent a corresponding plurality of data recording surfaces. Concentric tracks are defined on each of the recording surfaces and are provided with a physical track address determined by servo data written to each track.\nLogical track addresses are assigned to the tracks by positioning a first head adjacent a selected location of the data recording surface corresponding to the first head, measuring a head offset value for each of the remaining heads, and assigning logical track addresses to the tracks on each data recording surface in relation to the measured head offset values.\nIn this way, each logical cylinder will tend to include tracks with different physical track addresses, but the tracks in each logical cylinder will nonetheless be nominally aligned with the heads. This reduces the need for corrective seeks after head switching within the same logical cylinder, thus improving data transfer performance of the drive.\nThese and various other features and advantages which characterize the claimed invention will be apparent from a reading of the following detailed description and a review of the associated drawings."}
{"text": "Sutures are commonly used in medical procedures to sew tissue together in order to close tissue openings, cuts, or incisions during or after the medical procedure. Sutures are also used, in conjunction with anchors or other similar devices, to achieve or maintain traction, or other positioning of tissues or organs, during medical procedures. According to any of these uses and others, sutures are typically looped through the tissue and the one or more free ends knotted to maintain a desired position or tension of the suture. More specifically, a clinician may manually tie together a suture pair or knot a free end of a single suture to secure the appropriate positioning. According to embodiments where the suture is used to maintain traction, tension in the suture between the anchor and the knot tied adjacent the patient's skin ultimately provides the desired traction.\nAlthough knotting sutures may prove effective for certain procedures, there are a number of disadvantages of knotting sutures to secure tissues to one another and/or maintain a desired tension. For example, knot tying may require a considerable amount of time and may require a certain degree of dexterity. Further, knots may permanently fix a suture in place and, thus, may not be removed or adjusted once in place without removing the entire suture. Thus, although an additional knot may be created to increase suture tension, decreasing suture tension may require abandoning the previous suture and using another suture that may be knotted at the appropriate tension.\nU.S. Pat. No. 7,806,910 to Anderson teaches a suture clip comprising a plurality of flexible elements positioned together in a row. Specifically, first ends of the elements, which are bonded together at second ends thereof, are movable about living hinges and configured such that a tool may be used to urge the first ends of at least a portion of the flexible elements inward to define a slot between the first ends of the elements. A suture may be received within the slot, while the tool is actuating the flexible elements, and may later be held by the clip when the flexible elements are moved apart such that a tortuous path through the first ends of the elements is defined. The suture clip taught by Anderson, which appears to require a tool for actuation, is particularly suited for internal suturing, as described in the disclosure.\nThe present disclosure is directed toward one or more of the problems or issues set forth above."}
{"text": "The present invention resides generally in the field of regulation of myeloid cells such as stem cells or progenitor cells. More particularly, the present invention relates to the suppression of the proliferation of or the myeloprotection of mammalian myeloid cells using selected chemokines.\nAs further background, accessory cell-derived cytokines regulate proliferation/differentiation of hematopoietic stem and progenitor cells in vitro and in vivo. See, Broxmeyer, et al., \"The Production of Myeloid Blood Cells and Their Regulation During Health and Disease\", CRC Crit, Rev. Oncol/Hematol., Vol. 8, p. 173 (1988); and Broxmeyer, \"Biomolecule-cell Interactions and the Regulation of Myelopoiesis, An Update\", in Murphy Jr. (ed): Concise Reviews in Clinical and Experimental Hematology, Dayton, Ohio, Alpha Med Press, p. 119 (1992). Cytokines can have stimulating, enhancing, and/or suppressing activities mediated either directly on stem/progenitor cells and/or indirectly on accessory cells.\nA number of cytokines have been implicated in negative regulation. See, Broxmeyer, et al., \"The Production of Myeloid Blood Cells and Their Regulation During Health and Disease\", supra; Broxmeyer, \"Biomolecule-cell Interactions and the Regulation of Myelopoiesis, An Update\", supra; and Broxmeyer, \"Suppressor Cytokines and Regulation on Myelopoiesis: Biology and Possible Clinical Uses\", Amer. J. Ped. Hematol/Oncol, Vol. 14, p. 22 (1992). Suppression can be mediated by biological molecules termed cytokines, such as macrophage inflammatory protein (MIP)-1.alpha., a heparin binding protein originally identified by its capacity to cause a localized inflammatory reaction after injection into the footpads of C3H/HeJ mice. See, Wolpe, et al., \"Macrophages Secrete A Novel Heparin-Binding Protein with Inflammatory and Neutrophil Chemokinetic Properties\", J. Exp. Med., Vol. 167, p. 570 (1980); Davatelis, et al., \"Cloning and Characterization of a CDNA for Murine Macrophage Inflammatory Protein (MIP), a Novel Monokine with Inflammatory and Chemokinetic Properties\", J. Exp. Med., Vol. 167, p. 1939 (1988); and Sherry, et al., \"Resolution of the Two Components of Macrophage Inflammatory Protein 1, and Cloning and Characterization of One of Those Components, Macrophage Inflammatory Protein 1B\", J. Exp. Med., Vol. 168, p. 2251 (1988).\nMIP-1.alpha., but not a closely related family member MIP-1.beta., suppressed proliferation of a subset of murine (mu) stem cells (day 12 colony forming-unit-spleen (CFU-S)) ex vivo and mu colony forming unit-A (an apparently early progenitor cell) (see, Graham et al., \"Identification and Characterization of an Inhibitor of Hematopoietic Stem Cell Proliferation\", Nature, Vol. 334, p. 442 (1990)), as well as mu and human (hu) growth-factor stimulated multipotential (CFU-GEMM), erythroid (BFU-E) and granulocyte macrophage (CFU-GM) progenitor cells (see, Broxmeyer, et al., \"Enhancing and Suppressing Effects of Recombinant Murine Macrophage Inflammatory Proteins on Colony Formation in vitro by Bone Marrow Myeloid Progenitor Cells\", Blood, Vol. 76, p. 1110 (1990)) in vitro. Suppressive effects of MIP-1.alpha. were apparent on more immature populations of progenitors which were stimulated to proliferate by a combination of two or more early acting growth stimulating cytokines. See, Broxmeyer, et al., \"Enhancing and Suppressing Effects of Recombinant Murine Macrophage Inflammatory Proteins on Colony Formation in vitro by Bone Marrow Myeloid Progenitor Cells\", Blood, Vol. 76, p. 1110 (1990); and Bodine, et al., \"Effects of Hematopoietic Growth Factors on the Survival of Primitive Stem Cells in Liquid Suspension Culture\", Blood, Vol. 78, p. 914 (1991). These effects appeared to be directly on the progenitors themselves. See, Broxmeyer, et al., \"Macrophage Inflammatory Protein (MIP)-1.beta. Abrogates the Capacity of MIP-1.alpha. to Suppress Myeloid Progenitor Cell Growth\", J. Immunol., Vol. 147, p. 2586 (1991). MIP-1.alpha. did not suppress proliferation of the more mature progenitors which were stimulated to proliferate by a single cytokine. See, Broxmeyer, et al., \"Enhancing and Suppressing Effects of Recombinant Murine Macrophage Inflammatory Proteins on Colony Formation in vitro by Bone Marrow Myeloid Progenitor Cells\", supra; and Broxmeyer, et al., \"Macrophage Inflammatory Protein (MIP)-1.beta. Abrogates the Capacity of MIP-1.alpha. to Suppress Myeloid Progenitor Cell Growth\", supra.\nMIP-1.alpha. has recently been shown to have in vivo suppressive effects on cycling rates of CFU-S, CFU-GEMM, BFU-E an CFU-GM when administered to mice (see, Maze, et al., \"Myelosuppressive Effects in vivo of Purified Recombinant Murine Macrophage Inflammatory Protein-1 Alpha\", J. Immunol., Vol. 149, p. 1004 (1992); Dunlop, et al., \"Demonstration of Stem Cell Inhibition and Myeloprotective Effects of SCI/rhMIP-1.alpha. in vivo\", Blood, Vol. 79, p. 2221 (1992); and Lord, et al., \"Macrophage-Inflammatory Protein Protects Multipotent Hamatopoietic Cells from the Cytotoxic Effects of Hydroxyurea in vivo\", Blood, Vol. 79, p. 2605 (1992)), and in this context was myeloprotective for the drugs cytosine arabinoside (see, Dunlop, et al., supra) and hydroxyurea (see, Lord, et al., supra.) MIP-1.beta. was not myelosuppressive in vivo. See, Maze, et al., supra.\nMIP-1.alpha. also inhibited proliferation of an unstimulated cytotoxic T-cell line, CTLL-R8 (see, Oh, et al., Identification of Cell Surface Receptors for Murine Macrophage Inflammatory Protein-1.alpha.\", J. Immunol., Vol. 147, p. 2978 (1991)), and modulated macrophage function, including induction of the release of tumor necrosis-factor, interleukin (IL)-1.alpha. and IL-6 in vitro. See, Fahey, et al., \"Macrophage Inflammatory Protein 1 Modulates Macrophage Function\", J. Immunol., Vol. 148, p. 2764 (1992). MIP-1.beta., when present in excess compared to MIP-1.alpha. in vitro, blocked both the suppressive effects of MIP-1.alpha. on myeloid progenitors (see, Broxmeyer, et al., \"Macrophage Inflammatory Protein (MIP)-1.beta. Abrogates the Capacity of MIP-1.alpha. to Suppress Myeloid Progenitor Cell Growth\", supra), and the cytokine-inducing effects of MIP-1.alpha. on macrophages. See, Fahey, et al., supra.\nMIP-1.alpha. and MIP-1.beta. are members of a larger family of molecules variously termed small inducible proteins, intercrine cytokines (see, Wolpe, et al., \"Macrophage Inflammatory Proteins 1 and 2: Members of a Novel Superfamily of Cytokines\", FASEB, J., Vol. 3, p. 2565 (1989); Oppenheim, et al., \"Properties of the Novel Proinflammatory Supergene `Intercrine` Cytokine Family\", Ann. Ref. Immunol., Vol. 9, p. 617 (1991); and Schall, \"Biology of the Rantes/Sis Cytokine Family\", Cytokine, Vol. 3, pp 165 (1991)), and more recently, chemokines. They are linked by amino acid homology, chromosome location and the presence in their primary sequence of 4 position invariant cysteine residues. The hu MIP-1 family, located on chromosome 17 and having a c--c motif, includes MIP-1.alpha. (=LD78), MIP-1.beta. (=Act 2), Macrophage Chemotactic and Activating Factor (MCAF=muJE) and RANTES. The hu MIP-2 family located on chromosome 4 and having a c-x-c motif, includes GRO-.alpha. (also called melanoma growth stimulating factor=muKC), MIP-2.alpha. (=GRO-.beta.), MIP-2.beta. (=GRO-.gamma.), Platelet Factor 4 (PF4), IL-8 (=neutrophil activating peptide (NAP)-1) and NAP-2 (which derives from platelet basic protein an its derivatives connective tissue activating peptide III and .beta.-thromboglobulin)."}
{"text": "1. Technical Field\nThe present invention relates to technology for adjusting image formation conditions of an image forming device.\n2. Related Art\nIn image forming devices such as copiers or printers, recording methods such as electrophotography, thermal printing, inkjet printing or the like are used. Image forming devices have been subjected to various improvements in accordance with the characteristics of these recording methods, in order to fulfill such user demands as high-speed, and high-quality image formation.\nFor example, image forming devices employing electrophotography utilize the electrostatic effect of toner, which is made of microscopic particles, and therefore the image quality tends to drop due to fluctuations in the image darkness that occur when the toner is affected by environmental or temporal changes in temperature or humidity, or due to tiny variations in the structural components of the image forming device. Accordingly, electrophotographic image forming devices are commonly provided with a function for adjusting image darkness, in order to sustain a stable image quality.\nImage darkness adjustment means the forming of a standard image, such as a pattern image or patch image or the like, on a recording material (paper or the like) or image-carrying member (photosensitive member, intermediate transfer belt or the like), based on a target darkness that has been stored beforehand, and if the darkness of the formed standard image does not match the target darkness, adjusting the conditions relating to the image formation (hereinafter referred to as “image formation conditions”), such as the exposure potential during exposure or the amount of toner that is discharged during development or the like, such that the darkness of the formed standard image becomes the target darkness. By performing such an image darkness adjustment, images having a constant quality can be formed over a long period of time.\nGenerally, there are two kinds of image darkness adjustment, namely “automatic darkness adjustment” in which the image forming device itself performs the image darkness adjustment automatically when the power is turned on, when the apparatus is idle, or at predetermined time intervals; and “user calibration” in which the user instructs the image forming device to perform an image darkness adjustment, for example, when the user feels that the image quality has dropped. In these two kinds of image darkness adjustment, the specific methods for adjusting the image darkness may be the same, or they may be different.\nFurther, in addition to these two kinds of image darkness adjustment, image forming devices are normally also provided with a function by which the user can operate an operation section or the like of the image forming device to change the image formation conditions (hereinafter referred to as “manual darkness adjustment”). This manual darkness adjustment is often performed when the formed image does not have a level of image darkness or a color hue that is desired by the user. That is to say, the above-described automatic darkness adjustment and user calibration are adjustments for attaining a standard image quality, whereas the manual darkness adjustment is enables a user to select an image quality.\nThe automatic darkness adjustment and the user calibration are both adjustments for attaining a standard image quality. Therefore, after either an automatic darkness adjustment or a user calibration has been carried out, the image formation conditions should be the same, and if an identical image is output, then the image quality should be the same.\nHowever, if the method for adjusting the image darkness differs between the automatic darkness adjustment and the user calibration, then the image formation conditions do not necessarily match after the respective image darkness adjustments have been carried out, so there is a risk that the output results will be different. A technology has been developed with an object of addressing this problem by generating correction data from print output characteristic information, that corresponds respectively to a “device calibration”, which is equivalent to the automatic darkness adjustment, and a “software calibration”, which is equivalent to the user calibration.\nHowever, in these conventional technologies, the result of carrying out both the automatic darkness adjustment and the manual darkness adjustment, or both the user calibration and the manual darkness adjustment, is not taken into consideration.\nAs an example, a case will be explained where the image formation conditions are adjusted to a user's liking; that is, the image formation conditions after a manual darkness adjustment has been carried out are “lighter” than the image formation conditions after an automatic darkness adjustment has been carried out, that is, the standard image formation conditions. When the manual darkness adjustment is carried out, and a first automatic darkness adjustment is started in a state in which the image formation conditions are “lighter” than the standard image formation conditions, then the image forming device will form a standard image that is “lighter” than standard. Thus, a correction is performed to form a “darker” image, for example a correction is performed by which the look-up table is overwritten in a manner resulting in a larger output for a given input signal (see FIG. 12). After this automatic darkness adjustment, the image formation conditions approach the standard conditions, but in this situation, the user's preferences are not reflected, and the user again has to perform a manual darkness adjustment in order to set image formation conditions that are “lighter” than the standard image formation conditions. Thus, when the automatic darkness adjustment is performed for a second time, a correction is performed such that an image is formed which is even “darker” than the previous image. When such an adjustment is carried out repeatedly, the manual darkness adjustment and the automatic darkness adjustment will make repeated corrections in opposite directions. As a result, the second correction will be larger than the first and the third will be larger than the second, and the correction amount will continue to increase.\nFurthermore, when a large correction is performed as described above, the image forming device is not able to form an image using the original number of gradations, resulting in a problem that, for example, tone jumps brought about by the reduction of the number of gradations occur, and the reproducibility of the original image is compromised."}
{"text": "Conventionally, as heat-resisting steels for exhaust valves of automobile engine valves, there have widely been used 21-4N steel (JIS specification: SUH35), that is, a high Mn heat-resisting steel, and improved steels thereof which are good high temperature strength and oxidation resistance, and low cost.\nFor the face part of engine valves, high wear resistance is required because of continual contact with a valve seat. Accordingly, for the face part of valves using the 21-4N steel or improved steels thereof, usually, the built-up of Stellite etc. is done to thereby reinforce the hardness and wear resistance at high temperatures.\nMoreover, as a valve material used for portions exposed to a higher load, there is used in part a precipitation strengthening-type heat-resisting alloy including a lot of Ni and having an enhanced high temperature strength by precipitating γ′ (gamma prim) being an intermetallic compound, or NCF751 being a super heat-resisting alloy. However, since these alloys contain a lot of Ni, there is such a problem of increasing the cost.\nHowever, as the result of the tightening of environmental regulations in recent years, the efficiency and power of gasoline engines are enhanced to raise the combustion temperature, and therefore, a request is placed for a heat-resisting steel for valves which is low cost and excellent in high temperature strength as compared with the above-described heat-resisting alloys.\nIn order to answer the request, Japanese Patent Application Laid-Open No. 2001-323323 (Patent Document 1) proposes a production method of an engine valve, in which a base material formed by adding appropriately Mo, Nb and V besides C, N, Mn, Ni and Cr to a base of inexpensive Fe-based heat-resisting steel, and suppressing as much as possible the use of expensive raw materials such as Ni is used, which is subjected to a solution heat treatment at 1100 to 1180° C. and, after that, is subjected to forging in a temperature range of 700 to 1000° C. to form a valve having been subjected to an aging treatment of accumulating residual strain by machining intended for strain age hardening, thereby enhancing the hardness of the face part of the engine valve to 400 HV or more and controlling the overaging and softening even in the use in high temperature regions.\nFurthermore, Japanese Patent Application Laid-Open No. 2002-294411 (Patent Document 2) and Japanese Patent Application Laid-Open No. 3-177543 (Patent Document 3) propose engine valve materials obtained by adding, as an improved material of 21-4N steel being a high Mn heat-resisting steel, alloying elements such as Mo, W, Nb and V to thereby promote solid solution strengthening or precipitation strengthening and to improve high temperature strength and wear resistance."}
{"text": "The present invention relates to the well known process (hereafter “Process”) for the cryogenic separation of an air feed wherein:\n(a) the air feed is compressed, cleaned of impurities that will freeze out at cryogenic temperatures such as water and carbon dioxide, and subsequently fed into an cryogenic air separation unit (hereafter “cryogenic ASU”) comprising a main heat exchanger and a distillation column system which are contained in a large insulated box (generally referred to as the “cold box” in the industry);\n(b) the air feed is cooled in the main heat exchanger by indirectly heat exchanging the air feed against at least a portion of the effluent streams from the distillation column system;\n(c) the cooled air feed is separated in the distillation column system into effluent streams including a stream enriched in nitrogen, a stream enriched in oxygen and, optionally, respective streams enriched in the remaining components of the air feed including argon, krypton and xenon; and\n(d) the distillation column system typically comprises a first column (hereafter “high pressure column” or “HP column”) which separates the air feed into effluent streams including a nitrogen-enriched vapor stream and a crude liquid oxygen stream; and a second column (hereafter, “low pressure column” or “LP column”) which (i) operates at a relatively lower pressure than the HP column, (ii) separates the crude liquid oxygen stream into effluent streams including an oxygen product stream and one or more additional nitrogen-enriched vapor streams and (iii) is thermally linked with the HP column such that at least a portion of the nitrogen-enriched vapor from the HP column is condensed in a reboiler/condenser against boiling oxygen-rich liquid that collects in the bottom (or sump) of the LP column.\nMore specifically, the present invention relates to the known embodiment of the Process wherein the refrigeration extracted from liquefied natural gas (hereafter “LNG”) is utilized in order to provide the refrigeration necessary when at least a portion of the oxygen product is desired as liquid oxygen. In particular, the refrigeration is extracted from the LNG by indirectly heat exchanging the LNG in a heat exchanger against one or more nitrogen-enriched vapor streams withdrawn from the distillation column in order to liquefy such nitrogen-enriched stream(s). The skilled practitioner will appreciate the contrast between using LNG to liquefy such nitrogen-enriched stream(s) and the more conventional way of providing the refrigeration necessary to make liquid oxygen product. In particular, the more conventional way consists of turbo expanding a working fluid (typically either nitrogen or air).\nA key to the present invention is what happens to the nitrogen-enriched stream(s) that are liquefied against the boiling LNG. In particular, whereas the prior art introduces such stream(s) into the distillation column system, the present invention introduces such stream(s) into a heat exchanger (preferably the main heat exchanger) to be indirectly heat exchanged against at least a portion of the air feed to the distillation column system in order to liquefy at least a portion of the air feed to the distillation column system. In other words, whereas the prior art provides the LNG-derived refrigeration directly to the distillation column system, the present invention provides such refrigeration to the air feed. As further discussed herein, this has the advantage of both reducing the vapor feed to the high pressure column (thereby by allowing a smaller HP column at a smaller capital cost) and avoiding a safety hazard when, as per the prior art, the liquefied nitrogen is introduced into the distillation column directly after being indirectly heat exchanged against natural gas. In particular, in the event there is a defect in the heat exchanger used for the natural gas/nitrogen heat exchange such that natural gas leaks into the nitrogen, the leaked natural gas will be introduced directly into the distillation and thus have the potential to form very hazardous mixtures with oxygen.\nThe above described safety hazard is an important consideration because it leads to some of the unique features found in the below described prior art processes that utilize the refrigeration contained in LNG to aid in liquefaction.\nGB patent application 1,376,678 (hereafter “GB '678”) teaches the very basic concept of how LNG refrigeration may be used to liquefy a nitrogen stream. The LNG is first pumped to the desired delivery pressure then directed to a heat exchanger. The warm nitrogen gas is cooled in said heat exchanger then compressed in several stages. After each stage of compression, the now warmer nitrogen is returned to the heat exchanger and cooled again. After the final stage of compression the nitrogen is cooled then reduced in pressure across a valve and liquid is produced. When the stream is reduced in pressure, some vapor is generated which is recycled to the appropriate stage of compression.\nGB '678 teaches many important fundamental principles. First, the LNG is not sufficiently cold to liquefy a low-pressure nitrogen gas. In fact, if the LNG were to be vaporized at atmospheric pressure, the boiling temperature would be typically above −260° F., and the nitrogen would need to be compressed to at least 15.5 bara in order to condense. If the LNG vaporization pressure is increased, so too will the required nitrogen pressure be increased. Therefore, multiple stages of nitrogen compression are required, and LNG can be used to provide cooling for the compressor intercooler and aftercooler. Second, because the LNG temperature is relatively warm compared to the normal boiling point of nitrogen (which is approximately −320° F.), flash gas is generated when the liquefied nitrogen is reduced in pressure. This flash gas must be recycled and recompressed.\nU.S. Pat. No. 3,886,758 (hereafter “U.S. '758”) discloses a method wherein a nitrogen gas stream is compressed to a pressure of about 15 bara then cooled and condensed by heat exchange against vaporizing LNG. The nitrogen gas stream originates from the top of the lower pressure column of a double-column cycle or from the top of the sole column of a single-column cycle. Some of the condensed liquid nitrogen, which was produced by heat exchange with vaporizing LNG, is returned to the top of the distillation column that produced the gaseous nitrogen. The refrigeration that is supplied by the liquid nitrogen is transformed in the distillation column to produce the oxygen product as a liquid. The portion of condensed liquid nitrogen that is not returned to the distillation column is directed to storage as product liquid nitrogen.\nEP 0,304,355 (hereafter “EP '355”) teaches the use of an inert gas recycle such as nitrogen or argon to act as a medium to transfer refrigeration from the LNG to the air separation plant. In this scheme, the high pressure inert gas stream is liquefied against vaporizing LNG then used to cool medium pressure streams from the air separation unit (ASU). One of the ASU streams, after cooling, is cold compressed, liquefied and returned to the ASU as refrigerant. The motivation here is to maintain the streams in the same heat exchanger as the LNG at a higher pressure than the LNG. This is done to assure that LNG cannot leak into the nitrogen streams, i.e. to ensure that methane cannot be transported into the ASU with the liquefied return nitrogen. The authors also assert that the bulk of the refrigeration needed for the ASU is blown as reflux liquid into a rectifying column.\nU.S. Pat. Nos. 5,137,558, 5,139,547, and 5,141,543 (hereafter “U.S. '558”, “U.S. '547”, and “U.S. '543” respectively) provide a good survey of the prior art up to 1990. These three documents also teach the state-of-the-art at that time. U.S. '558 teaches cold compression to greater than 21 bara such that the nitrogen pressure exceeds the LNG pressure. U.S. '547 deals with the liquefier portion of the process—key features are cold compression to 24 bara and refrigeration recovery from flash gas. U.S. '543 further teaches to use turbo-expansion in addition to LNG for refrigeration to liquefy nitrogen.\nThere is little new art in the literature since the early 90's because the majority of applications for recovery of refrigeration from LNG (LNG receiving terminals) were filled and new terminals were not commonly being built. Recently, there has been resurgence in interest in new LNG receiving terminals and therefore the potential to recover refrigeration from LNG.\nWith respect to the ASU operation, a fundamental teaching of U.S. '758 is illustrated in FIG. 1. The facility includes an LNG-based nitrogen liquefier (2) and a cryogenic ASU (1). In this example, the cryogenic ASU includes a higher pressure column (114), lower pressure column (116), and main exchanger (110). Feed air 100 is compressed in 102 and cleaned of impurities that will freeze out at cryogenic temperatures such as water and carbon dioxide in unit 104 to produce stream 108. Stream 108 is cooled in 110 against returning gaseous product streams, to produce cooled air feed 112. Stream 112 is distilled in the double column system to produce liquid oxygen 158, high pressure nitrogen gas (stream 174) and low pressure nitrogen gas (stream 180). The nitrogen gases 174 and 180 are warmed in the main exchanger 110 to produce streams 176 and 182. Streams 176 and 182 are processed in the LNG-based nitrogen liquefier to create liquefied nitrogen product stream 184 and liquid nitrogen refrigerant stream 186. Liquid nitrogen refrigerant stream 186 is introduced into the distillation columns through valves 136 and 140.\nThe principle laid-out in FIG. 1 is also taught in JP 2005134036, JP 55-77680 (JP 1978150868), U.S. Pat. Nos. 4,192,662, 4,054,433, as well as the above referenced U.S. '758 and EP '355. There are two disadvantages related to processes based on FIG. 1. Firstly, should there be a leak of hydrocarbon into ASU refrigerant stream 186, that hydrocarbon will concentrate in the bottom of the lower pressure column and in liquid oxygen stream 158. Since build-up of hydrocarbon in oxygen is to be avoided, for reasons of safety, steps must be taken to ensure that such a leak does not occur in the LNG-based nitrogen liquefier. Secondly, since all of the incoming air to the cryogenic ASU (stream 108) is introduced to the higher pressure column as vapor, this requires a bigger diameter (and thus higher capital cost) for the higher pressure column.\nIt is therefore desired to provide an efficient process that transports refrigeration of the LNG-based nitrogen liquefier to the cryogenic ASU without the disadvantages associated with directly injecting potentially hydrocarbon laden liquid nitrogen to the distillation columns.\nAs used herein, “LNG-based nitrogen liquefier” shall be defined as a system that uses the refrigeration contained in LNG to convert gaseous nitrogen into liquid nitrogen. Typical of such systems, the nitrogen will be compressed in stages. If the compression is performed with a cold-inlet temperature, the LNG will be used to cool the compressor discharge by indirect heat exchange. Cooling and or liquefaction of the nitrogen will be accomplished, at least in part, by indirect heat exchange with warming or vaporizing LNG. Examples of LNG-Based Nitrogen Liquefiers can be found in the above referenced GB '678, U.S. '558, U.S. '547, and U.S. '543."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for recognizing a security code having an electromagnetic band gap (EBG) pattern.\n2. Description of the Related Art\nGenerally, Microwave Band Gap (MBG) structures or EBG structures are implemented on microstrips, and are used for various purposes, e.g., improving the performance of an antenna, improving the power efficiency of an amplifier, realizing the high quality (Q) of a resonator, controlling harmonic components, and designing a new type of duplexer. Microstrip circuits applied to such EBG structures can be manufactured using a method of punching a dielectric substrate, a method of etching a ground surface to have a uniform shape, and a method of changing a microstrip line itself."}
{"text": "Triple negative breast cancer (TNBC) represents a significant proportion (about 20-25%) of breast cancer patients. TNBC is characterized by the absence of HER2, estrogen receptor (ER), and progesterone receptor (PR). TNBC has a poor prognosis, and no targeted approach to a therapy has been found to date."}
{"text": "1. Field of the Invention\nThe present invention relates to a vehicular adaptor system which can adapt a portable telephone set to a vehicle having no vehicular telephone or to a vehicular telephone installed in a vehicle, used in a multiple carrier access system, such as a cellular mobile telephone system.\n2. Description of the Related Art\nRecently, portable telephones have become smaller in size and lighter in weight, i.e. convenient to carry. Accordingly, due to increased popularity, their numbers are increasing. A portable telephone can be considered a compact and light weight modification of the vehicular telephone due to its lower output power than the vehicular telephone. Since it is easily portable, it is used almost in the same way in the service area of the vehicular telephone. The portable telephone can be used in a vehicle; however, if it is used as it is, the shielding effect of the metal body of the vehicle decreases the radio frequency wave propagation, i.e. output power and receiving sensitivity, of the portable telephone, and may result in its operation being disabled. Therefore, a radio frequency booster/pre-amplifier is installed in the vehicle and a vehicular antenna is installed outside the vehicle.\nA prior art vehicular adaptor system for a portable telephone set is hereinafter described with reference to a perspective view in FIG. 1. A portable telephone set 11 is composed of a main body of the portable telephone set and an antenna (not seen there due to being integrated in the main body). The main body includes an earphone, a speaker, a key board, and a display for displaying a call number dialed onto the key board, etc. When the portable telephone set 11 is rested in a cradle 12' of a vehicular adaptor 12, an end A.sub.4 of cradle 12' and an end A.sub.2 of portable telephone set 11 are electrically contacted so as to communicate control signals, audio signals and an adaptor signal informing the fact that the portable telephone set 11 is mounted on the adaptor, between the adaptor 12 and the portable telephone set 11. The antenna integrated in the main body of the portable telephone set 11 is coupled with a radio frequency coupler (not seen in the figure) integrated in a side A.sub.3 of cradle 12' so as to face the integrated antenna. Adaptor 12 communicates control signals of booster 13 and radio frequency signals to/from booster/pre-amplifier 13 by means of a control signal cable and a radio frequency cable. Adaptor 12 further communicates up-link/down-link data and audio signals to/from a remote control unit (referred to hereinafter as control unit) 14 by means of another cable.\nFunctions of the above-described adaptor system for the portable telephone are summarized as follows:\n1. The portable telephone set brought into the vehicle is rested on vehicular adaptor 12. Then, a battery installed in the portable telephone set is charged by a charging circuit of the adaptor.\n2. 0.6 watt output power of the portable telephone set is fed via vehicular adaptor 12 to radio frequency booster 13, wherein 0.6 watt is amplified up to 3 watt.\n3. A radio frequency signal to be input to portable telephone set 11 is amplified so as to compensate for coupling losses at the radio frequency coupler, etc., and accordingly becomes the level as strong as a level of the portable telephone set when used outside the vehicle.\n4. Control unit 14 is used in the same way as a vehicular telephone set, i.e. is used as a hand set for talking, for dialing a call number via its key board and for displaying the phone number, etc. thereon.\nTo use portable telephone 11 in a vehicle in which a vehicle telephone has been already installed, all of FIG. 1 adaptor system is independently required in addition, except for a commonly used vehicular antenna.\nThe problems of the prior art FIG. 1 adaptor system include the following:\n(1) Having to control both of the control unit and the radio frequency power booster/pre-amplifier, the adaptor unit becomes complicated in configuration and large in size.\n(2) For connecting the integrated antenna of the portable telephone set to the radio frequency booster/pre-amplifier, a radio frequency coupler is required in the adaptor, causing complicated structure and precise mechanism.\n(3) The vehicular adaptor becomes large in size caused the system to become large and occupy a large space in the vehicle.\n(4) A portable telephone set to be used in a vehicle already having a vehicle telephone installed thereon requires a considerable investment for the duplicated radio frequency circuits, as well as a large space therefor."}
{"text": "The construction industry has experienced a growing trend in the use of insulated concrete forms (ICFs), wherein forms for pouring concrete are constructed from multiple modular form units. Each unit includes inner and outer sidewalls, at least one of which is formed of foamed polystyrene, foamed polyurethane, or other cellular plastics or insulating materials. The sidewalls of the form units are stacked or otherwise interconnected at the construction site to form opposing insulated inner and outer form walls between which concrete is poured. The insulated form walls are then left with the poured concrete at the site to define a portion of the poured concrete wall(s) of the structure being constructed, resulting in concrete walls with insulated surfaces. Examples of insulated concrete forms and form units of this nature can be found, for example, in U.S. Pat. Nos. 4,706,429 and 4,866,891 to Young; U.S. Pat. Nos. 4,765,109 and 4,889,310 to Boeshart; U.S. Pat. Nos. 5,390,459 and 5,809,727 to Mensen; and U.S. Pat. No. 6,314,697 to Moore.\nAs these patents illustrate, it is common to have each sidewall of a form unit bear tongue-and-groove structures (or other interfitting structures) at its edges so that the inner sidewall of each form unit can be interfit at its edges to inner sidewalls of other form units, thereby allowing the inner sidewalls to be combined to form an inner wall of a concrete form. The outer sidewalls can likewise include interfitting structure allowing them to be combined into an outer form wall. Additionally, the inner and/or outer sidewalls often include “webs,” structures which are generally formed of plastic and which extend within and engage the foamed insulating material of the sidewalls. Connecting members which are often referred to as “ties” or spacers then extend between the inner and outer sidewalls and engage their webs to hold the sidewalls in opposing parallel relationship. When the concrete is poured between the sidewalls to solidify, the ties are left embedded within the concrete and maintain the insulated sidewalls as cladding on the opposing sides of the concrete wall.\nWhile form units and forms of the foregoing nature are beneficial in that they conveniently use the forms for casting the concrete walls as insulating cladding for the walls, and they eliminate any need to disassemble or remove the forms after the walls are poured, they suffer from the disadvantage that their form units—being formed of a pair of sidewalls (generally foamed of bulky foamed plastic) joined by spacers—occupy substantial volume, and are therefore expensive to ship. Some of the aforementioned patents address this disadvantage by providing detachable/reattachable spacers which rigidly but disconnectably affix the sidewalls together. Such form units allow users to provide sidewalls and spacers separately, whereby the sidewalls of each form unit are stacked and shipped separate from the spacers (and thus without including a wasted intermediate space between the sidewalls), and each form unit can then be assembled at the construction site by fastening the spacers between the sidewalls. However, these forms trade shipping costs for labor costs, since hundreds or even thousands of spacers must be installed between the sidewalls to construct the form units and forms.\nTo overcome the foregoing difficulties, some ICF manufacturers have developed concrete form units wherein the spacers are pivotally affixed to their opposing sidewalls, with the various spacers thereby effectively form parallelogram linkages with the sidewalls. As a result, the sidewalls can be brought together (their intermediate space may be eliminated) by moving the sidewalls in opposing longitudinal directions. Examples of such arrangements are found in U.S. Pat. No. 3,985,329 to Liedgens, and U.S. Pat. Nos. 6,230,462 and 6,401,419 to Beliveau. Form units of this nature are useful because the concrete form units may be collapsed (their sidewalls may be brought into closely spaced relationship with the intermediate space eliminated), and the form units may be stacked in close relationship for shipping. The form units may then be readily unloaded at the construction site, unfolded to their expanded states, and assembled to construct larger concrete forms. However, these are disadvantageous in that the parallelogram linkage arrangement gives rise to “racking”: the sidewalls, when collapsed, are offset and do not rest end-to-end, and therefore generate unused volume which is effectively wasted during shipping. This is undesirable since the form units are already quite bulky, and expensive to ship. Additionally, while users need not install the spacers between the sidewalls because the spacers are already pivotally affixed therebetween, the expanded form units are subject to buckling because the spacers do not rigidly situate the sidewalls in spaced relation. Such buckling can lead to difficulties, particularly when using the concrete form units to construct a larger concrete form."}
{"text": "This present invention relates to a garment with a magnetic vent closure which enables a garment wearing individual to facilitate air ventilation through a vent window thereof with convenience by application of a magnetic element thereto.\nAs an increasing number of people evaluate outdoor activities, various function-focused garments are on the market. When on hiking, for example, the hiker needs to take off his/her jacket from time to time for cooling off the heated body. In particular, motorcycle riders suffer inconvenience due to a venting jacket according to the prior art which employs a zipper. Since a motorcycle rider typically wears gloves, it has been difficult to open the zipper during the motorcycle riding for air ventilation in the prior art.\nAccordingly it is understood that air ventilation is in strong demand for garments which appeal to those who enjoy outdoor activities such as motorcycle riding and bicycle riding."}
{"text": "Large boilers, for example, those used in conjunction with steam turbines for electric power generation, are fired by fuels such as coal, oil, gas or liquor. Supporting igniter burners are typically associated with each of the main burners. Because the igniter burners are typically fired with relatively expensive fuels, they are operated only intermittently. More particularly, the igniter burners are preferably fired only upon initial start up of the boiler and thereafter they are only selectably fired for short intervals to light off or support flame at the particular main burner(s) associated with the igniter burner(s).\nThe prior art has evolved a variety of flame detection techniques for monitoring boiler fires to detect the presence or absence of flame in the boiler regions supported by the various igniter burners. If flames are extinguished in a particular region of the boiler, then the \"no flame\" condition must be quickly identified or else the main burners continue to supply fuel which may potentially explode if it is not evenly and continuously ignited. Accordingly, highly reliable flame monitoring techniques are required for continuously detecting the presence of flame at regions within the boiler adjacent to each of the burners which fire the boiler.\nThe apparatus to be described in this application is suitable for use with two types of boiler/burner configurations; namely, \"wall\" (or \"opposed\") fired boilers, and \"corner\" (or \"vortex\") fired boilers. \"Wall\" or \"opposed\" fired boilers incorporate a series of burners mounted on two opposing walls of the four vertical walls of the boiler. Sighting tubes (pipes about 5 cm. in diameter) are positioned across the boiler walls (which are typically about 1.5 meters thick) beside and nearly parallel to each burner head. The sighting tubes are pointed approximately toward the expected location of burner flame. Flame detection apparatus is positioned to \"sight\" through each tube into the boiler region in which flame is expected.\n\"Corner\" or \"vortex\" fired boilers incorporate vertically separated stacks of burners in each of the four corners of the boiler. The flames produced by the burners merge in a central vortex within the boiler. The burners may be individually tilted in the vertical plane in order to better control the combustion characteristics and location of the fireball within the boiler. Sighting tubes for corner fired boilers must be flexible so that the flame detection apparatus can continuously track the flame as the burners tilt.\nSeveral prior art flame detectors examine the light emitted by the flame and, from the time variation characteristics of these emissions, determine whether a flame is located near to the burner (\"near flame\"); or, a fireball is present in the background (\"far flame\"); or, there is no detectable flame. By monitoring flame flicker (i.e. time variations in the light signal emitted in the frequency band(s) under consideration) such prior art detectors attempt to derive a binary signal representative of \"flame\" and \"no flame\" conditions. Pre-determined factors such as the geometry of the detector, the wavelength band it is capable of examining, and the frequency band being monitored affect the characteristics of flame flicker and correspondingly determine the ability of such detectors to accurately detect the presence or absence of flame under varying conditions.\nThe best prior art flame detectors for use on opposed fired boilers appear to be those which utilize two separate linear arrays of detectors aligned horizontally and vertically to facilitate \"X-Y\" scanning of selected sub-regions within a region where flame is expected, through electronic selection of an appropriate detector pair. Typically, a zero-crossing waveform shaping analysis is performed on the electronic signals produced by each of the two selected detectors, to generate two bi-level output signals. The output signals are then correlated with one another (prior art detectors of this sort do not however perform true signal correlation, because they work only with binary (i.e. two level) approximations of the detector output signals, rather than with the direct analog outputs of the detectors). If the two signals are highly similar to one another then the correlation result approaches unity. Normally, a result which exceeds some predetermined threshold is accepted as indicating the presence of flame. If the two signals are highly dis-similar to one another then the correlation result approaches zero. A result which does not exceed the aforementioned threshold is normally taken to indicate a \"no flame\" condition. In some cases, automatic tracking techniques are employed to locate points of maximum correlation in an effort to minimize generation of false \"no flame\" alarms. It will thus be understood that the prior art is susceptible to error, in that the cumulative approximations inherent in the operation of prior art detectors may result in a \"no flame\" alarm when flame is in fact present; or, may indicate that flame is present when no flame is in fact present. The prior art tends to overcompensate for these contingencies by allocating flame determination thresholds which minimize generation of false \"flame present\" signals. However, this necessarily decreases the ability of such prior art devices to respond to flame conditions having light emission characteristics which encompass a large dynamic range.\nThe inventors believe that superior results may be obtained by concentrating on factors other than flame flicker. More particularly, the inventors believe that superior results may be obtained by analyzing the time.fwdarw.frequency spectral characteristics of light emitted in different visual and infra-red wavebands from the region in which flame is expected, and comparing those characteristics with prestored spectral signatures representative of flame. The present invention accordingly compares short term estimates of the visible and infra-red auto-spectra, the infra-red to visible transfer function, and the infra-red to visible coherence (all of which are hereinafter defined and explained in greater detail), with prestored signatures characteristic of \"flame\" and \"no flame\" conditions. The auto-spectra, transfer function and coherence function are used to characterize the relationship between two signals in selected frequency bands. It is this relationship or pattern which is used to identify the flame."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the fields of genetics, protein biochemistry, and oncology. More particularly, it concerns the use of exosomal genomic DNA and proteins in genetic analysis and treatment.\n2. Description of Related Art\nPancreaticoduodenectomy (Whipple procedure) can be curative for PDAC patients if tumors are detected early with clear surgical margins. Due to the late diagnosis of pancreatic cancer, only around 15% of patients present with surgically resectable tumors (Conlon et al., 1996). Studies comparing stage of disease with outcome following surgery suggest that death rates for PDAC would be reduced if the disease were diagnosed at an earlier stage (Bilimoria et al., 2007).\nIn addition to direct cell-to-cell contact via soluble factors, such as cytokines and chemokines, there is emerging evidence that exosomes play a pivotal role in intercellular communication (Kahlert and Kalluri, 2013). Exosomes are small, membrane-bound vesicles with a size of 40-150 nm (Pan et al., 1985; Trams et al., 1981). They are secreted by many different cell types, such as cancer cells, mesenchymal cells, thrombocytes (Kahlert and Kalluri, 2013; Heijnen et al., 1999; Raposo et al., 1996), immune cells (Thery et al., 2009), platelets (Janowska-Wieczorek et al., 2005), and endothelial cells (Hergenreider et al., 2012). The first step in exosomes biogenesis involves the inward budding from the limiting membrane of late endosomes (Trajkovic et al., 2008). During this process, exosomes are packed with RNA molecules and proteins from the parental cell (Trams et al., 1981; Trajkovic et al., 2008). After the release into the extracellular space, tumor-derived exosomes can transfer proteins and RNAs with oncogenic activity to recipient cells (Kacharzewska et al., 2012; Grange et al., 2011; Peinado et al., 2012). Because exosomes are very stable under different conditions, they can protect their biological cargo against degradation and denaturation in the extracellular environment (Taylor and Gercel-Taylor, 2008). Genomic DNA in circulation is mainly contained in exosomes (Kahlert et al., 2014). Exosomes from astrocytes and glioblastoma cells carry mitochondrial DNA (Guescini et al., 2010). Furthermore, it has been shown that exosomes from glioblastoma cell lines contain small amounts of single-stranded DNA as well as high levels of transposable elements (Balaj et al., 2011).\nExosomes are found in all body fluids of cancer patients, such as serum, saliva, cerebrospinal fluid, bone marrow aspirates, eye exudate/tears, and ascites (Peinado et al., 2012; Lau et al., 2013; Choi et al., 2011). As such, exosomes are promising diagnostic and predictive biomarkers in cancer. However, genetic profiling studies on circulating DNA from cancer patients are confounded by the fact that the isolated DNA represents all cells of the body, thus making mutation and genetic defects challenging (Murtaza et al., 2013; Yong, 2014; Kirk, 2013; Corwley et al., 2013).\nSeveral exosomes markers have been proposed and include members of the tetraspanin family (CD9, CD63, CD81), members of the endosomal sorting complexes required for transport (ESCRT; TSG101, Alix), and heat shock proteins (Hsp60, Hsp70, Hsp90) (Taylor and Gercel-Taylor, 2011). Epithelial tumor cells secrete exosomes carrying the epithelial cell adhesion molecule (EpCAM) (Taylor and Gercel-Taylor, 2008; Silva et al., 2012; Runz et al., 2007). Melanoma-derived exosomes contain the tumor-associated antigen Mart-1 and tyrosinase-related protein-2 (TYRP2) (Peinado et al., 2012; Mears et al., 2004; Andre et al., 2002). Exosomes from gastric cancer, breast cancer, and pancreatic cancer carry members of the human epidermal growth factor receptor (HER) family (Adamczyk et al., 2011; Baran et al., 2010; Ciravolo et al., 2012). However, none of these markers are specific to cancer-derived exosomes and specific isolation of exosomes from the serum of cancer patients remains a challenge due to the lack of specific markers that can be used to identify and distinguish cancer exosomes from exosomes produced by other cells. A marker for cancer-derived exosomes will significantly increase the sensitivity of detection for low frequency mutations in circulation. Thus, a procedure to specifically detect and isolate cancer cell-derived exosomes in circulation is needed."}
{"text": "The present invention relates to a system for checking the function of a catalyst used to purify the exhaust gas in an internal combustion engine. In particular, the invention is concerned with a function diagnostic system for an exhaust gas purifying apparatus in an internal combustion engine, which diagnostic system has suitable for evaluating the exhaust gas purifying performance for a specific component.\nHeretofore there has been adopted a method in which O.sub.2 sensors are disposed upstream and downstream of a catalyst, and whether the catalyst, is deteriorated or not is judged by comparing the outputs of both sensors. In this case, it is actually the O.sub.2 storage capacity out of the various catalyst functions that is measured, rather than the conversion efficiency of each exhaust gas component.\nAccording to the above conventional method, when the conversion efficiency is determined independently of the O.sub.2 storage capacity, it is impossible to make an accurate detection of the conversion efficiency. For solving this problem, the use of a sensor for detecting a specific exhaust gas component, instead of an O.sub.2 sensor, may be effective in measuring the conversion efficiency of that exhaust gas component. An example of using an HC sensor for detecting the conversion efficiency of HC (hydrocarbon) is disclosed in Japanese Patent Laid Open No. 109021/92. According to the method disclosed therein, however, it is necessary to dispose HC sensors both upstream and downstream of a catalyst. Therefore, the environmental conditions, particularly around the upstream-side HC sensor, are severe (for example, high temperature and high exhaust gas component concentrations), thus resulting in the sensitivity and output characteristic of the sensor itself being deteriorated. This eventually exerts a bad influence on the diagnostic accuracy for catalyst deterioration."}
{"text": "The invention relates to a method of controlling the torque developed at a belt shaft of a belt retractor coupled to an electric motor. DE 100 18 972 discloses a belt retractor whose belt shaft is driven by an electric motor rather than by a coil spring. The electric motor is coupled to the belt shaft by a toothed belt. Assuming the electric motor to be supplied with a constant amount of current, friction and the flexing work of the toothed belt result in a hysteresis of the seat belt forces, being different when the seat belt is wound up and when it is unwound.\nIt is an object of the invention to essentially avoid hysteresis forces in order to keep the pressure exerted by the seat belt on the body of the vehicle occupant essentially constant.\nIn the method according to the invention, torque developed at a belt shaft of a belt retractor coupled to an electric motor is controlled by detecting a direction of rotation of the belt shaft and, depending on the direction of rotation as detected, changing electric current supplied to the electric motor from a first value to a second value that is different from the first value by an amount required to compensate for hysteresis.\nAccordingly, to equalize the seat belt forces the amount of electric current supplied to the electric motor is adjusted in such a way that it is increased in the winding direction and decreased in the unwinding direction, resulting in substantially equal torque at the belt shaft and hence essentially equal seat belt forces when the seat belt is wound up and unwound.\nIn accordance with a preferred embodiment of the invention, a standstill of the belt shaft is also detected. The seat belt forces needed to keep the pull-out length of the seat belt constant are smaller than the forces needed to unwind it. Accordingly, to compensate for the smaller forces, a higher amount of electric current is supplied to the electric motor.\nIn case of an emergency, when strong deceleration forces act on the vehicle, inertia forces tend to move the body of the vehicle occupant away from the seat, entraining and unwinding thereby the seat belt at a speed that is higher than during normal movement of the vehicle occupant. In order to protect the vehicle occupant, such a rapid unwinding of the seat belt has to be avoided.\nIn accordance with another preferred embodiment of the invention, rotational speed of the belt shaft is measured in an unwinding direction and when a predetermined first threshold value of rotational speed is exceeded, the resistance against pulling out the seat belt is increased by increasing the amount of electric current supplied to the electric motor, the rate of increase being constant. Only when a second threshold value of rotational speed, corresponding to a threshold value of the amount of electric current, is reached, a mechanical locking mechanism on the belt retractor is activated. Otherwise a restriction of the seat belt is achieved by means of the electric motor only.\nA preferred use of the inventive method is in a retractor of the type comprising a frame, a belt shaft rotatably fixed in the frame and an electric motor coupled to the belt shaft by a drive belt, wherein the direction of rotation can be detected and, depending on the direction of rotation as detected, the motor current is adjusted."}
{"text": "Some aspects of the present invention relates generally to tools used to implant and remove neural interface devices within nervous systems. More particularly, some embodiments of the invention relates to device systems that can be used to insert and retract a range of microscale devices dependent upon the desired research or clinical purpose."}
{"text": "1. Field of the Invention\nThe present invention relates to a rotation correcting apparatus and an optical disk apparatus that is used to record or reproduce data from an optical disk.\n2. Related Art\nIn recent years, along the increase in rotation speed such as a CD-R/RW or a DVD-Drive, a drive unit that rotates a disk at near a limit rotation number of a disk motor and a pickup mechanism is being developed progressively. For example, techniques of increasing the data reading performance of a wobbling disk or an eccentric disk during a high-speed rotation and improving the accuracy of track jump and layer jump in a DVD double-layer disk is proposed (see Japanese Patent Application Laid-open Publication No. 2003-263760).\nAccording to the technique described in the above document, a signal that follows the eccentricity or wobbling of the optical disk is A/D converted, and is stored in the memory. Based on the data stored in this memory, a pickup or a feed motor that shifts the pickup to a radial direction of the optical disk is controlled.\nAccording to the above document, both a tracking servo system and a feed motor control system are subjected to eccentricity correction. However, the output from a feed motor control circuit that controls the feed motor makes no change irrespective presence or absence of eccentricity correction. A phase of a control signal from the feed motor control system is delayed from that of a control signal of the tracking servo system due to a phase delay in the feed motor control circuit. Therefore, when eccentricity correction is carried out in both the tracking servo system and the feed motor control system, the operation of the feed motor becomes unstable, because the actual eccentricity of the optical disk and the operation of the feed motor are not in the same phase. As a result, the servo control of the tracking servo system is not convergent."}
{"text": "Tumor growth and metastasis impacts a large number of people each year. In fact, it is estimated that well over 600,000 new cases of cancer will be diagnosed in the coming year in the United States alone (Varner, J. A., Brooks, P. C., and Cheresh, D. A. (1995) Cell Adh. Commun. 3, 367-374). Numerous studies have suggested that the growth of all solid tumors requires new blood vessel growth for continued expansion of the tumors beyond a minimal size (Varner et al. 1995; Blood, C. H. and Zetter, B. R. (1990) Biochim. Biophys. Acta. 1032:89-118; Weidner, N. et al. (1992) J. Natl. Cancer Inst. 84:1875-1887; Weidner, N. et al. (1991). N. Engl. J. Med. 324:1-7; Brooks, P. C. et al. (1995) J. Clin. Invest. 96:1815-1822; Brooks, P. C. et al. (1994) Cell 79:1157-1164; Brooks, P. C. et al. (1996). Cell 85, 683-693; Brooks, P. C. et al. (1998) Cell 92:391-400. A wide variety of other human diseases also are characterized by unregulated blood vessel development, including ocular diseases such as macular degeneration and diabetic retinopathy. In addition, numerous inflammatory diseases also are associated with uncontrolled neovascularization such as arthritis and psoriasis (Varner et al. 1995).\nNew blood vessels develop from pre-existing vessels by a physiological process known as angiogenesis (Varner et al. 1995; Blood and Zetter 1990; Weidner et al. 1992). This complex process requires cooperation of a variety of molecules including growth factors, cell adhesion receptors, matrix degrading enzymes and extracellular matrix components (Varner et al. 1995; Blood and Zetter 1990; Weidner et al. 1992). Thus, therapies designed to block angiogenesis may affect the growth of solid tumors. In fact, clear evidence has been provided that blocking tumor neovascularization can inhibit tumor growth in various animal models, and human clinical data is beginning to support this contention as well (Varner, J. A., Brooks, P. C., and Cheresh, D. A. (1995) Cell Adh. Commun. 3, 367-374).\nIt has also been proposed that inhibition of angiogenesis can be effected by (1) inhibition of release of “angiogenic molecules” such as βFGF (fibroblast growth factor), (2) neutralization of angiogenic molecules, such as by use of anti-βFGF antibodies, and (3) inhibition of endothelial cell response to angiogenic stimuli. This latter strategy has received attention, and Folkman et al., Cancer Biology, 3:89-96 (1992), have described several endothelial cell response inhibitors, including collagenase inhibitors, basement membrane turnover inhibitors, angiostatic steroids, fungal-derived angiogenesis inhibitors, platelet factor 4, thrombospondin, arthritis drugs such as D-penicillamine and gold thiomalate, vitamin D3 analogs, alpha-interferon, and the like that might be used to inhibit angiogenesis. For additional proposed inhibitors of angiogenesis, see Blood and Zetter 1990; Moses et al. (1990) Science 248:1408-1410; Ingber et al. (1988) Lab. Invest., 59:44-51; and U.S. Pat. Nos. 5,092,885, 5,112,946, 5,192,744, and 5,202,352.\nTo block angiogenesis, many investigators have also focused on growth factors and cytokines that initiate angiogenesis (Varner et al. 1995; Blood and Zetter 1990; Weidner et al. 1992; Weidner et al. 1991; Brooks et al. 1995; Brooks et al. 1994; Brooks et al. 1997). However, there is a large number of distinct growth factors and cytokines which have the capacity to stimulate angiogenesis. The therapeutic benefit of blocking a single cytokine may have only limited benefit due to this redundancy. Accordingly, what is needed is other anti-angiogenic targets for inhibiting angiogenesis and proteolysis."}
{"text": "Microprocessor designers have increasingly endeavored to improve performance in various microprocessors by defining an instruction which does a common function in order to save executing a number of instructions to perform the same function. For this reason, some microprocessors have an instruction for setting a selected bit field in an operand to all 1's, or for clearing a selected bit field in the operand to all 0's. Typically, a sequence of microcode instructions is used to implement such a bit field instruction by shifting 0's or 1's, as appropriate, into the operand at a designated location. As a result, several iterative cycles are required to execute a bit field instruction.\nAn object of the present invention is to overcome the performance delay caused by iterative execution of microcode for a bit field instruction."}
{"text": "A content sharing platform may receive millions of messages from users desiring to share media content such as audio, images, and video between user devices (e.g., mobile devices, personal computers, etc.). The media content of these messages may be associated with a common geolocation, a common time period, a common event, and so forth."}
{"text": "The present invention relates generally to the field of sensors, and more particularly to an improved pixel architecture with increased signal to noise ratio and with threshold response.\nTwo-dimensional large-area sensor arrays have wide ranging applications in medical imaging, optical scanning, chemical or radiation detection, temperature sensing, and other arts. Such sensors comprise a grid of pixels, with each pixel composed of at least a sensor device and one or more transistors comprising a circuit for converting sensory data (such as optical images, x-ray images with or without the assistance of a scintillation layer, chemical presence, temperature, etc.) to electrical signals (typically data used for processing, storing, driving a display or an alarm, etc.) The sensor is generally a reverse biased diode, and the transistor(s) is generally a thin-film field effect transistor (TFT). Conveniently, each of the components of the pixel may be formed of hydrogenated amorphous silicon (a-Si) to take advantage of the ease of fabrication of large area arrays, low leakage current devices, and other benefits of using this amorphous material. (A component will be referred to herein as being formed of a-Si when at least one layer of the component is formed of said material. For example, as used herein, an a-Si TFT refers to a thin film transistor having an amorphous silicon channel layer.)\nIt should be clear from the above that the sensors of the type relevant to the present invention are designed to sense one or more of a variety of sensory data (e.g., light, radiation, chemical, thermal, etc.). An important subset of this group of sensors is comprised of sensors optimized to sense image data. While the discussions herein are broadly applicable to sensory data, the following discussion will focus on image data for clarity of explanation.\nIt is known that a critical parameter of an array of image sensors is the signal to noise ratio of the array. Namely, a reduction of noise as compared to signal is generally desirable. Signal amplifiers are typically employed to boost sensed signals for downstream use. It is common practice to connect an entire row of image sensors to a single amplifier. There are numerous advantages to this architecture, including minimizing the number of contacts to the array, maximizing the fill factor (defined as the fraction of the area of the array sensitive to the incident image), etc.\nHowever, the noise in such an architecture is related to both the readout noise due to the large capacitance of the long data line, and the pixel noise which is a function of the pixel architecture itself and which is amplified together with the signal itself. For example, the thermal noise of the readout amplifier is proportional to the capacitance on the input of the amplifier.\nAccordingly, it has been suggested that each pixel be provided with its own amplifier circuitry. In copending U.S. patent application Ser. No. 08/699,875, to Street et al. (referred to herein as \"Street et al, now U.S. Pat. No. 5,831,258.\"), a pixel circuit and array incorporating such a circuit is disclosed. Typically, such a circuit consists of a photosensor with inherent capacitance for storing a charge under reversed bias conditions. The stored charge is depleted by light incidence. Determining the intensity of the incident light is accomplished by examining the remaining charge in the photosensor. The photosensor may be connected in reverse bias mode between a data line and a bias line. The photosensor is connected to a transistor or transistors such that the voltage through the transistor(s), between the bias line and the data line, is a function of the amount of light incident upon the photosensor. This structure is referred to as an amplified pixel.\nIn addition to increasing the signal to noise ratio, it is also a perennial goal in the image sensor art to provide increased resolution from sensor arrays. For example, a pixel having an area on the order of (150 .mu.m).sup.2, as taught by Street et al., corresponds to a resolution of approximately 170 dpi, and qualifies as a low resolution sensor by current measures.\nIt is generally understood that reducing pixel area is a means of increasing resolution. However, there are several penalties generally paid when pixel area is reduced. The first is that the area of the photosensor is reduced. This results in a decrease in sensitivity of the photosensor. The second is that one or more of the dimensions of the pixel transistors are reduced. In the case of an amplified pixel, this runs the risk of reducing gain below an acceptable threshold. For example, in the case of a single transistor amplifier, an a-Si TFT with a ratio of channel width (W) to channel length (L) of 4 provides a desired gain of approximately 10. A reduction in the width to length ratio (W/L), however, results in a reduction in gain. The combination of current process limitations leading to a minimum device dimension of L.gtoreq.5 .mu.m for a-Si devices, combined with the low carrier mobility of a-Si, means that a channel of 5 .mu.m.times.20 .mu.m is as small as current technology will permit while providing adequate gain. However, there is a desire to overcome this barrier and produce smaller pixels for increased resolution.\nIt will be appreciated that there are two parameters of a sensor circuit at odds with the desire to reduce pixel size. The first is the desire to maximize the sensitivity of the pixel. This can be accomplished by increasing the physical size of the photosensor part of the circuit, but at the cost of increasing the entire pixel size. Rather than increasing pixel size, it is common in the art to increase the area of the cell dedicated to the photosensor, referred to as the fill factor. As between two sensor pixels who differ only in photosensor area, the one with the larger photosensor area (greater fill factor) will generally provide better sensitivity. However, for a given pixel size, generally the larger the area occupied by transistor(s), the smaller the area available for the photosensor. Given the lower limit on transistor size, a tradeoff must therefore be made between resolution and sensitivity.\nThe second parameter at odds with the desire to reduce pixel size is the complexity of the circuit. A minimal circuit requires a photosensor and a selector switch (e.g., a transistor). However, this circuit does not address the need to do any amplification or in situ signal processing, etc. While the added circuitry is desirable from a pixel performance perspective, it undesirably consumes pixel area (and reduces yield). A tradeoff must thus be made between pixel size or fill factor and circuit complexity."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming system that includes an image forming apparatus mounted with a plurality of components and a lifetime management device that manages information on the lifetime of the components.\n2. Description of the Related Art\nWhen a failure occurs in a part of an image forming apparatus such as a copier, a facsimile machine, and a printer, depending on the type of the part, the image forming apparatus cannot be used until the part is replaced by a new one, and this imposes inconvenience on a user.\nTherefore, various types of image forming systems have been proposed, which estimate the end of the service life of respective parts or components of an image forming apparatus to issue a request to replace a part which reaches its service life. For example, Japanese Patent Application Laid-open No. H9-146423 discloses such an image forming system. In the conventional image forming system, the number of prints produced by an image forming apparatus is calculated with respect to each of parts mounted thereon. When the number of prints exceeds a predetermined threshold for one of the parts, it is determined that the part has deteriorated to about the end of its service life due to printing operation. Then, a signal is transmitted to a remote monitoring device to request replacement of the part. In response to the request, a replacement worker, such as a service person, replaces the part before a failure occurs, which reduces the downtime of the image forming apparatus.\nSome of the various parts mounted on the image forming apparatus deteriorate depending on an operation amount specific to the parts rather than the number of prints. For example, a developing roller, which does not closely contact recording paper during the printing operation, deteriorates depending on the operation amount including the total rotation time and the total rotation distance rather than the number of prints (hereinafter, this type of part is referred to as “operation amount-dependent deteriorated part”). On the other hand, a cleaning member that removes paper dust from the recording paper and toner remaining on an intermediate transfer belt, and a fixing roller that strongly contacts the recording paper deteriorates depending on the number of prints rather than the operation amount (hereinafter, this type of part is referred to as “print volume-dependent deteriorated part”). If the number of prints is completely proportional to the operation amount, the end of service life of each part can be accurately estimated based on the number of prints. However, the number of prints and the operation amount do not show a proportional correlation.\nA factor of this relates to an idle operation at the start and the end of a print job. Specifically, the printing operation of the image forming apparatus includes a single printing operation in which an image is formed only on one recording sheet, and a continuous printing operation in which images are continuously formed on a plurality of recording sheets. In either printing operation, an idle operation, in which a sheet is not fed to respective parts, is performed when a job starts and ends. The idle operation is performed for the same time period in the single printing operation and the continuous printing operation. Accordingly, in the single printing operation, the percentage of the idle operation time in the total operation time is high compared to the continuous printing operation. Further, in the continuous printing operation, as the number of continuous prints decreases, the percentage of the idle operation time increases. Therefore, when a user frequently performs single printing and rarely performs continuous printing of relatively a few copies, the operation amount is considerably large, while the number of prints is relatively small. In such a case, if the service life is estimated based only on the number of prints, the operation amount-dependent deteriorated part may reach the end of its service life before the estimated one. On the contrary, when a user performs continuous printing of a large number of copies frequently and single printing relatively infrequently, the operation amount is considerably small, while the number of prints is relatively large. In such a case, if the service life is estimated based only on the number of prints, it is likely to be determined that the operation amount-dependent deteriorated part is almost at the end of its service life even if sufficient time remains until the end.\nTherefore, the present inventors have studied to estimate service life based on the operation record of each part, instead of the number of prints, for the operation amount-dependent deteriorated parts. However, the number of prints more or less relates to the progress of deterioration even in the operation amount-dependent deteriorated parts according to the type of recording paper. Specifically, paper dust can adhere to the operation amount-dependent deteriorated parts such as the developing roller, which does not closely contact recording paper, via the intermediate transfer belt and a photosensitive drum, which closely contact recording paper. When a user mainly uses recording paper likely to generate paper dust, such as low-quality paper, the amount of adhering dust considerably increases. As a result, the number of prints relates relatively closely to the progress of deterioration even in the operation amount-dependent deteriorated parts. In such a case, if the service life of the operation amount-dependent deteriorated part is estimated based only on the operation record, estimation accuracy decreases."}
{"text": "The invention relates to an optical multi-gate element with an acousto-optical modulator (AOM), which exhibits at one side at least two optical waveguide connections, in particular monomode optical waveguides. The end surfaces of the waveguides are disposed in the focal plane of a lens situated between the optical waveguides and the acousto-optical modulator.\nSuch an arrangement is known from \"Journal of Lightwave Technology\", Vol. Lt-2, No. 2, pages 108 to 115. In this case, a gradient lens is coaxially associated with each individual optical waveguide forming a port connection. A transformation of the field widths of the Gaussian beams is caused by the lenses. In consequence of this, the field width of the beams on the path is increased by the AOM as compared with the field width in the optical waveguides.\nIn the case when the AOM is not excited acoustically, the rays are not diffracted and further transmitted on a direct path. In the case of acoustic excitation, beams impinging at the Bragg angle to the acoustic wavefront are deflected by twice the Bragg angle. On account of this effect, an optical switch can be constructed with an AOM, as in the known case. Since the frequency of the deflected beam is also additionally increased or reduced by the acoustic frequency, an AOM multi-port element is also suitable for the formation of a local oscillator beam (LO) for optical heterodyne reception. The combined use of an AOM as switch and for the formation of an LO beam is known from the DE-OS No. 3,506,884 in an application for a heterodyne OTDR.\nAn OTDR (Optical time domain reflectometer) is a device with which the attenuation behaviour of an optical waveguide can be measured from one end, whereby the components of an optical transmission signal (laser pulse) which are backscattered by the longitudinal positions of the optical waveguide are evaluated. In order to achieve an OTDR with a large range, its optical components should have little attenuation, so that the intensity of the back-scattered signals, which decreases with increasing distance from the OTDR, is still sufficiently large to ensure an adequate signal-to-noise ratio.\nA multi-port element of the initially mentioned type consists of numerous individual elements, which necessarily result in further losses. Besides the unavoidable material-dependent transmission attenuations, there are insertion attenuations, which are dependent upon the quality and geometric accuracy of the coupling-in of the optical beams. Thus, a careful relative adjustment of the individual components is essential for a low-attenuation multi-port element of the initially mentioned type. Furthermore, long free beam paths through air should be avoided, since, on the one hand, a small overall length of the multi-port element is, of course, to be aimed at, and since, on the other hand, irradiation losses, in particular due to beam broadening, unavoidably arise on the free beam paths.\nWithout particular measures, the direct beam and the beam deflected in the case of acoustic excitation can be detected separately only at a relatively large distance from the AOM by the lenses which cannot be constructed arbitrarily small, because the Bragg angle is very small. In order to avoid an excessively large distance, a deflecting prism is provided, in the known case mentioned initially which prism, however, causes additional attenuation losses and requires a high degree of expenditure on adjustment."}
{"text": "Poly(para-phenylene) (PPP) is a fully aromatic, rigid rod polymer with unique structural and conductive properties. As an engineering plastic, its attractiveness arises from its thermal stability (mp&gt;500.degree. C.), high strength, chemical inertness, and solvent resistance. When doped with either n- or p-type dopants, the polymer forms highly conducting charge transfer complexes with conductivities up to 500 S/cm. However, the structural properties which make PPP so attractive also make it a difficult polymer to synthesize. In addition, many of the observed properties of the polymer depend on the method of production.\nPrevious methods of producing PPP directly have met with only limited success. For example, oxidative cationic polymerization of benzene to produce PPP has been attempted. However, only short oligomers of ten to fifteen repeat units containing mixtures of linear 1,4- and non-linear 1,2-units were formed.\nPolymerizations using nickel catalyzed aryl coupling of 1,4-dihalobenzenes were attempted. While this method produced a completely linear molecule, only short oligomers consisting of ten to twelve units were formed.\nThe problem with these direct synthetic methods is that the inherent insolubility of the polymer causes it to precipitate out of solution before high molecular weight materials can be formed. Electrochemical coupling of benzene has also been used, but the resulting film is insoluble and composed of a mixture of 1,4- and 1,2- units.\nIn order to circumvent the problem of the inherent insolubility of PPP in production and processing, soluble precursor methodologies have been developed. For example, polymers of 1,3-cyclohexadiene (CHD) have been used as a soluble precursor polymer. In particular, poly(cyclohexadiene) has been reacted with bromine and then pyrolyzed to eliminate HBr. Unfortunately, this polymerization route also produces a precursor polymer with a mixture of 1,4-and 1,2-linkages. In addition, the elimination reaction is not very efficient, since HBr readily reacts with unsaturated intermediates.\nRecently, the efficient production of PPP has been reported, via the pyrolysis of a soluble precursor polymer prepared from the radical polymerization of the acetyl and methoxycarbonyl derivatives of 5,6-dihydroxy-1,3-cyclohexadiene (DHCD); see, e.g., D. G. Ballard et al., Macromolecules, Vol. 21, pp. 294-304 (1988) and D. R. McKean, Macromolecules, Vol. 20, pp. 1787-1792 (1987). The starting cis-diol is produced by the microbial oxidation of benzene. The precursor films are soluble and can be processed before pyrolyzing into the final polymer. However, the radical polymerization produces about 85% 1,4-units and 15% 1,2-units. The 1,2-units create \"kinks\" in the polymer, thereby reducing the elimination efficiency of the precursor and the mechanical properties of the final polymer.\nIn order for the good mechanical properties of PPP to be realized, an aspect ratio of at least 100 consecutive linear units per 1,2-unit must be obtained. All previous routes to PPP have either produced low molecular weight materials due to insolubility of the growing polymer, or have incorporated a significant amount of 1,2-linkages in the chains, or both. Hence, it is desirable to find an exclusively 1,4-polymerization method which can be used in combination with the efficient precursor method described above."}
{"text": "Exemplary embodiments disclosed herein pertain to current limiters, such as circuit breakers. More particularly, exemplary embodiments disclosed herein pertain to current limiters in an environment which may induce transient surges.\nFIG. 1 is a schematic diagram wherein a prior art telecommunications power card 2 is depicted. The purpose of telecommunications power card 2 is to provide constant voltage to a load that is connected to it.\nOn the exterior of the card, common node 4 and negative VIN node 6 are disposed to connect to a voltage source providing, for example, −48 volts. The term “node” refers to any equipotential point in an electrical device, either at a terminus, or not at a terminus. Voltage sources of −48 volts are widely used in telecommunications applications; nearly all telephone central offices are powered from −48 volt sources wherein the “hot” side of the supply is negative with respect to ground.\nHot swap controller 8 is disposed on prior art telecommunications power card 2, and is coupled to positive ground node 4 and negative VIN node 6 from which it derives approximately −48 volts when connected to an active power source providing −48 volts. Hot swapping (also known as hot plugging) is the ability to add, remove and replace components of a system while the system is powered. In the context of telecommunications applications, the hot swapping of power sources, such as batteries or power supplies, is commonly performed.\nCapacitor 10 is electrically coupled to hot swap controller 8 at nodes 12 and 14. Increases of voltage correspondingly store excess charge in the capacitor, thus moderating the change in output voltage and current induced by a transient power event.\nSince it is desirable to provide less than the −48V input for a load that is coupled to telecommunications power card 2, (e.g. +3.3 volts), a DC to DC converter 16 is provided on telecommunications power card 2 and is electrically coupled to hot swap controller 8 at nodes 12 and 14. Thus, capacitor 10 and DC to DC converter 16 are arranged in parallel fashion. DC to DC converter 16 is a load with respect to the power output of hot swap controller 8.\nIn electrical engineering, a DC to DC converter is a circuit which converts a source of direct current from one voltage to another. A DC to DC converter may serve to isolate an input power source (i.e. hot swap controller 8 from the load connected to the output of the DC to DC converter at common node 18 and VOUT node 20. DC to DC converters are commercially available as blocks or modules from Datel, Power Trends, Integrated Power Designs, or others.\nFIG. 2 is a schematic diagram which depicts one form of a current limiter 22 of the prior art. Current limiter 22 is coupled to common node 4 and negative VIN node 6. These nodes provide power to the current limiter 22 which in turn provides power to the DC-DC converter 16 via VOUT node 24 and common node 4. Current source 26 is coupled to common node 4, and provides a constant reference current. It is comprised of voltage reference source 28 resistor R1, amplifier 30, and transistor 32 which are arranged in a well known configuration to provide constant current I1 based on the reference voltage provided by voltage reference source 28. Note that in FIG. 2, voltages are referred to −VIN at node 6.\nResistor R2 is coupled to current source 26 at node 34 and to negative VIN node 6. The constant current provided by current source 26 flows through resistor R2 producing reference voltage V1 across resistor R2 such that:V1=I1·R2\nResistor RSENSE is coupled to negative VIN node 6 and to pass transistor 36 which acts as a circuit breaker at node 38. When the circuit is not broken by pass transistor 36, current flows between VOUT node 24 and negative VIN node 6. The purpose of resistor RSENSE is to sense the current 12 flowing between negative VIN node 6 and VOUT node 24 and produce a voltage VSENSE at node 38 such that:VSENSE=I2·RSENSE \nComparator 40 is coupled to node 34 and to node 38 such that it compares voltage VSENSE to voltage V1. The output of comparator 40 is coupled to latch 42 indicating when the circuit should be broken or modified. During normal operation, VSENSE is less than V1, so comparator 40 signals that the circuit should not be broken. When VSENSE exceeds V1, comparator 40 signals that the circuit should be broken or modified in some way.\nConstant current I1 is only a small fraction of the current limit. R2 can be chosen to produce the reference voltage corresponding to the desired limit for current I2 as sensed by resistor RSENSE. Since the resistance value of R2 is multiplied by constant current I1, higher reference voltages can be achieved by increasing the resistance value of R2.\nLatch 42 is coupled to the output of comparator 40 and to pass transistor 36. When comparator 40 signals that the circuit should be broken, latch 42 propogates the signal to pass transistor 36, breaking the circuit. Latch 42 maintains the signal to break the circuit until receiving a reset signal via reset node 44. As will be appreciated by those skilled in the art, there are many alternatives to Latch 42 for various applications.\nFIG. 3 is a graph depicting a constant-current threshold current limiter of the prior art. In FIG. 3, the vertical axis indicates current level and the horizontal axis indicates time increasing from left to right. The dashed line labeled ICB represents a constant-current threshold in common use. When this threshold is exceeded, it is an indication that the circuit breaker must interrupt current flow between the power source and the load. The solid line in FIG. 3 depicts an example of a current level which is initially constant. At a certain point in time, the current level suddenly surges to a dramatically higher level, which is above the constant current limit.\nSince the current level exceeds the constant-current threshold, the circuit is broken. In some cases, this is undesirable. For example, if an additional power source is brought online, the voltage of the combined power sources may suddenly increase. Since a smoothing capacitor is typically coupled in parallel with the power supply and the load, the capacitor begins charging to the new voltage. The transient increase in current caused by the capacitor charging to the new voltage causes the circuit to break prematurely. Although the current surge may be relatively brief, and the current level decays exponentially as the capacitor approaches the new voltage, it can cause premature circuit breakage because the design of the prior art constant-current threshold circuit breaker is unforgiving.\nOne approach to this problem is to provide a current limiter which is temporarily disabled, for a fixed amount of time, to allow for a transient surge. This approach has the disadvantage that it does not protect the load from a damaging current surge such as that induced by a lightning strike or other dramatic power event. What is needed is a current limiter which is continuously online in order to protect the load, but which handles the transient surge caused by bringing a new or additional power source online; the prior art fails to properly address this problem.\nThese and other limitations of the prior art will become apparent to those of skill in the art upon a reading of the following descriptions and a study of the several figures of the drawing."}
{"text": "1. Field\nExemplary embodiments of the present invention relate to a pixel circuit and a display device using the pixel circuit. More particularly, exemplary embodiments of the present invention relate to a pixel circuit that can be driven at a high speed and display an image of uniform luminance, and a display device using the pixel circuit.\n2. Discussion of the Background\nVarious kinds of flat display devices developed in recent years are capable of reducing detriments of cathode ray tubes (CRT), such as their heavy weight and large size. Examples of flat panel displays are a liquid crystal display (LCD), a field emission display (FED), a plasma display panel (PDP), an organic light emitting diode (OLED) display, and the like.\nAmong the flat panel displays, the organic light emitting diode display which uses an organic light emitting diode (OLED), generates light by recombining electrons and holes, has a fast response speed, is driven with low power consumption, and has excellent emission efficiency, luminance, and viewing angle, such that it has recently been in the limelight.\nIn general, a plurality of pixels emitting light in the organic light emitting diode (OLED) display include an organic light emitting diode (OLED), and the organic light emitting diode (OLED) generates light of a predetermined luminance corresponding to a data current supplied from a pixel circuit.\nDigital driving, which is one method of expressing grayscale of the organic light emitting diode display, controls an on time of a pixel. When the digital driving is used in the organic light emitting diode, one frame is divided into a plurality of sub-frames, and a light emitting period of each sub-frame is set in order to display grayscale. Among a plurality of sub-frames constituting one frame, the pixel emits light during the selected sub-frame depending on an image signal for grayscale expression. That is, the sub-frame is selected according to the image signal, and is turned on to express the grayscale.\nThe above information disclosed in this Background section is only for enhancement of understanding of the background of the invention and therefore it may contain information that does not form the prior art that is already known in this country to a person of ordinary skill in the art."}
{"text": "This invention relates to semiconductor high voltage devices, and specifically to semiconductor high voltage devices with a voltage sustaining layer containing floating regions.\nIt is well-known that in many semiconductor devices, such as VD-MOST and SIT, a high sustaining voltage always accompanies a high specific on-resistance. This is due to the fact that, for a high sustaining voltage, thickness of a voltage sustaining layer should be large and doping concentration of the voltage sustaining layer should be low, so as the peak field does not exceed the critical field for breakdown −EC, which is normally expressed by EC=8.2×105×VB−0.2V/cm for silicon, where VB is the breakdown voltage of the voltage sustaining layer.\nIn a uniformly doped n-type voltage sustaining layer between p+-region and n+-region, in order to obtain a minimum specific on-resistance at a given breakdown voltage, a doping concentration ND and a thickness W of the voltage sustaining layer are optimized such that a maximum field is at p+-n-junction and its value is equal to EC, a minimum field is at n+-n-junction and equal to EC/3. For a silicon device,ND=1.9×1018×VB−1.4 cm−3  (1)W=1.8×10−2×VB−1.2ÿm−2  (2)(see, e.g., P. Rossel, Microelectron. Reliab., vol. 24, No. 2, pp 339-366, 1984).\nFor the VDMOST shown in FIG. 1A, a field profile in the voltage sustaining layer at VB is shown in FIG. 1B, where a slope of the field versus distance is qND/Es, Es is the permittivity of the semiconductor and q is the electron charge. The change of field through the n-region is qND/Es, 2EC/3. The relation between Ron and VB of a n-type voltage sustaining layer is then expressed byRon=W/qÿÿnND=0.83×10−8×VB2.5ÿ.cm2  (3)where ÿn is the mobility of the electron and ÿn=710×VB0.1 cm/V.sec is used for silicon.\nIn order to get even lower Ron at a given VB, some research has been done to optimize the doping profile instead of using a uniform doping, see: [1] C. Hu, IEEE Trans. Electron Devices, vol. ED-2, No. 3, p 243 (1979); [2] V. A. K. Temple et al., IEEE Trans. Electron Devices, vol. ED-27, No. 2, p 243 (1980); [3] X. B. Chen, C. Hu, IEEE Trans. Electron Devices, vol. ED-27, No. 6, p 985-987 (1982). However, the results show no significant improvement."}
{"text": "1. Field of the Invention\nThis invention relates to sign holder arrangements which are suspended to a ceiling rail of a store, and more particularly to a rotatable spool which has a stop feature thereon for permitting adjustment and control of the height and the alignment of a sign suspended from a ceiling.\n2. Prior Art\nThe use of signs in stores are more convenient when they are able to be supported from an overhead ceiling rail. Such signs permit shoppers to see advertisements about products which are easy to notice. To be able to do this well, a sign should be hung from the ceiling rail in an even manner. Since signs may come in various sizes, a support arrangement for those signs should be adjustable, so as to permit their level to be controlled and regulated, and their height as supported from ends of the sign, should be equivalent so as to not present an uneven or improperly supported sign.\nIt is an object of the present invention, to provide a sign holding mechanism which is an improvement over the prior art.\nIt is a further object of the present invention, to provide an attachable ceiling rail pod which permits accurate sign positioning when suspending a sign from such a ceiling rail.\nIt is a further object of the present invention, to provide a sign hanging system which is readily adjustable and accurately alignable with a second sign hanging arrangement.\nIt is a further object of the present invention, to provide a sign hanging system which may be easily attached to and removed from a ceiling rail.\nIt is also an object of the present invention, to provide such a sign supporting or suspending system which is economical to manufacture and very easy to use."}
{"text": "Every year millions of dollars are spent on advertising. However, companies are always looking for new ways to advertise their products or services. Furthermore, professional sports leagues use their advertising revenues to pay the players on the teams in the leagues and to pay for other operating expenses. For example, soccer leagues typically place advertising signs around the field and NASCAR® drivers place advertising on their vehicles and on their clothing. Like any other business, the sports leagues are always looking for new and improved ways to obtain advertising revenue.\nAccordingly, a new method of advertising is desired."}
{"text": "The present invention relates to a process for the extraction of cesium ions from aqueous solutions using an adduct compound comprising a macrocyclic polyether (crown ether) and an inorganic heteropoly acid component.\nThe processing and solidification of medium active aqueous wastes (MAW), developed during the reprocessing of irradiated nuclear fuels and/or nuclear breeder materials, e.g., in form of waste solutions, waste concentrates or waste sludges, is rendered more difficult by the presence of small quantities of the cesium isotopes Cs.sup.134 and Cs.sup.137. During the vitrification of MAW, cesium evaporates to a noticeable degree and is also easily leached out of the solidified products intended for permanent storage, such as glass blocks, cement blocks and bitumen blocks. Selective extraction of the Cs would considerably simplify the treatment of medium active wastes. Furthermore, if Sr.sup.90 were simultaneously extracted from the MAW, low radioactive aqueous wastes (LAW=low active waste) would be obtained since, after a relatively short decay period, practically the entire activity of MAW originates from the relatively long-lived radionuclides Cs.sup.137 (t.sub.1/2 =30 yrs) and Sr.sup.90 (t.sub.1/2 =26 yrs).\nAccording to prior-art processes, the extraction of Cs and Sr was achieved mainly by means of coprecipitation reactions. The coprecipitation reactions, however, did not result in satisfactory DF (decontamination factor) values for Sr and Cs. Therefore, different processes have been investigated which would make possible a selective extraction of the two radionuclides, either jointly or individually.\nThe extraction processes for Cs.sup.+ ions developed thus far, however, are not suited for the extraction of Cs.sup.+ from a typical MAW with a high concentration of NaNO.sub.3 and free nitric acid. The composition of such a typical MAW is shown in the following Table 1:\nTABLE 1 ______________________________________ Components Concentration (strongly acidic MAW) g/dm.sup.3 ______________________________________ Nitric acid 64 (1 mol/dm.sup.3) Sodium nitrate 42.5 (0.5 mol/dm.sup.3) Uranium 1.0 Lead 0.8 Calcium 0.2 Magnesium 0.2 Iron 0.15 Ruthenium 0.07 Potassium 0.05 Zinc 0.05 Copper 0.04 Manganese 0.02 Chromium 0.02 Nickel 0.01 Zirconium 0.005 Cesium 0.004 Cerium 0.004 Strontium 0.002 Antimony 0.002 Niobium 0.001 ______________________________________\nThe extraction processes for Cs.sup.+ ions that have been developed use mostly organic reagents which are added either to the organic phase or to the aqueous phase to be extracted. These reagents, however, would be destroyed either hydrolytically or by oxidation, or would be made ineffective by protonization if they came into contact with the MAW whose essential components are such as those described in Table 1 above.\nCs.sup.+ selective heteropoly acids are known which are very stable in the highly acid environment provided by an MAW having a composition such as described in Table 1. The Cs.sup.+ selective heteropoly acids are used as extraction agents in the form of an organic phase with nitrobenzene as a solvent. For example, 12-molybdophosphoric acid (HPMo) is a Cs.sup.+ selective heteropoly acid which can be dissolved in nitrobenzene. In this form, the heteropoly acids are equally unsuitable for the continuous Cs.sup.+ extraction from an MAW of the type described in the above Table 1, since they bleed considerably from the nitrobenzene into the aqueous phase due to their good water solubility. They would have to be replaced continuously.\nIf HPMo is used, the molybdenum would enter into the MAW in large quantities and would furthermore interfere with the subsequent processing of the MAW, such as during vitrification. Finally, depending on the Cs.sup.+ concentration of the MAW and the heteropoly acid concentration in the nitrobenzene phase, there is the danger of precipitate formation which could lead to considerable disruptions in a continuously operating system.\nIn contrast to the above prior-art water soluble extraction agents for Cs.sup.+, dibenzo crown ethers have a very low water solubility and are thus, at least in part, suitable for use in a continuous extraction process.\nSince crown ethers are neutral ligands, the anion in the aqueous solution being treated is also extracted during the extraction. This, however, significantly influences the extraction coefficient. In this way, simple cesium salts, such as chloride and nitrate salts, are very difficult to extract with the use of crown ethers in polar solvents, with the exception of nitromethane solutions and nitrobenzene solutions. For this reason, the prior art started to combine crown ethers with known Cs.sup.+ specific reagents which contain voluminous, polarizable anions. Up to now, this has been done by adding the Cs.sup.+ specific reagents, in most cases picric acid to the aqueous phase to be extracted,, and extracting with the aid of crown ethers, in most cases dibenzo-18-crown-6 (DB-18-C-6) or dicyclo-hexyl-18-C-6, in polar solvents, in most cases nitrobenzene.\nFor the extraction of Cs.sup.+ from nuclide containing solutions, J. Rais and P. Selucky proposed two extraction systems which use DB-18-C-6. In one method, as described in Czechoslovakian Pat. No. 149403, DB-18-C-6 was employed in an organic phase and dipicrylamine was added to the aqueous phase to be extracted. In a second method, as described in Czechoslovakian Pat. No. 149404, DB-18-C-6 was employed in an organic phase and sodium tetraphenylborate was added to the aqueous phase to be extracted.\nThese processes, however, are limited to alkaline Cs.sup.+ solutions having a pH of 11 to 13 inasmuch as dipicrylamine is protonized in the acid range, and sodium tetraphenylborate is hydrolyzed in the acid range. Moreover, both processes only function well in the absence of large quantities of Na.sup.+ and K.sup.+.\nFurthermore, the same authors, as disclosed in Czechoslovakian Pat. No. 165751, have produced an adduct of DB-18-C-6 with 12-molybdophosphoric acid which permits a column chromatographic extraction of the Cs.sup.+ from waste solutions. Large quantities of Na.sup.+ and K.sup.+, however, also present a problem here. The adduct of DB-18-C-6 with 12-molybdophosphoric acid was prepared by a two phase reaction in which an organic phase comprising a solution of the crown ether in dichloromethane was shaken with an aqueous phase comprising a solution of HPMo in 2 mol/1 HNO.sub.3. This patent also refers generally to adducts with 12-tungstophosphoric acid (HPW). A similar disclosure of an adduct of DB-18-C-6 with molybdophosphoric acid appears in J. Radioanal. Chem. Vol. 35 (1977), pages 351 to 359."}
{"text": "1. Field of the Invention\nThe present invention relates to a camera system, a digital camera, and a method for controlling a digital camera and, more particularly, to a camera system that makes backup data of image data obtained by photographing by a digital camera, a digital camera which can transfer image data obtained by photographing to the outside, and a method for controlling a digital camera.\n2. Description of the Related Art\nIn recent years, digital cameras such as digital still cameras or digital video cameras have become popular. Such cameras pick the image of an object by using a solid-state image pickup element such as a CCD (Charge Coupled Device) that records digital image data indicating the object image obtained by the image pickup on a mobile recording medium such as a smart medium or a compact flash.\nIn these kinds of digital cameras, as a method for protecting image data recorded on the recording media, the following two methods are employed.                1. Pieces of image data to be protected are selected by a user one by one to give a read-only attribute to each selected image data.        2. Read-only attributes are given to all pieces of recorded image data at once.        \nExamples of the protect methods of the two types will be described below in detail. With reference to FIGS. 15A to 15D, a method for giving a read-only attribute to each selected piece of image data will be described below.\nA user depresses a menu/execution switch (not shown) prepared for displaying a menu on an image monitor 40 constituted by a liquid crystal display or the like arranged on a digital camera. In accordance with this operation, a menu screen shown in FIG. 15A is displayed on the image monitor 40 of the digital camera.\nThe user selects a protect menu indicated by a key-shaped mark from the menu screen such that an operation switch (not shown) on which an arrow key indicating four moving directions, e.g., upper, lower, left, and right directions in a display area of the image monitor 40 is operated. In this manner, a protect menu shown in FIG. 15B is displayed on the image monitor 40 of the digital camera. Therefore, the user designates “1 frame set” of the protect menu by operating the operation switch. In accordance with this, on the image monitor 40 of the digital camera, an image indicated by image data recorded on a recording medium by photographing in advance is displayed as a thumbnail. For this reason, the user designates an image to be protected from the thumbnail by operating the operation switch. In accordance with the designation, a screen is displayed on the image monitor 40 of the digital camera. That screen urges the user to perform re-confirmation, as shown in FIG. 15C. When image data corresponding to a designated image is to be protected, the user depresses the menu/execution switch. When the image data are to be canceled, the user depresses a cancel/return switch (not shown).\nWhen the user depresses the menu/execution switch, a read-only attribute is set with respect to image data corresponding to an image designated by the user in the digital camera to protect the image data, and a key-shaped mark is displayed at the upper right portion of the display area of the image monitor 40 as shown in FIG. 15D.\nA method for giving read-only attributes to all pieces of recorded image data at once will be described below with reference to FIGS. 16A to 16C.\nFirst, the user depresses the menu/execution switch. In accordance with this, a menu screen shown in FIG. 16A is displayed on the image monitor 40 of the digital camera.\nThe user selects a protect menu from the menu screen by operating the operation switch in the same manner as described above. In this manner, a protect menu shown in FIG. 16B is displayed on the image monitor 40 of the digital camera. In this state, the user designates “all frame set” of the protect menu by operating the operation switch. In accordance with this, a screen, which urges the user to perform re-confirmation shown in FIG. 16C, is displayed on the image monitor 40 of the digital camera. When all the pieces of recorded image data are protected at once, the user depresses the menu/execution switch. When the protection is canceled, the user depresses the cancel/return switch.\nWhen the user depresses the menu/execution switch, in the digital camera, read-only attributes are set for all the pieces of recorded image data to protect the image data.\nEven though a read-only attribute is given to the image data recorded on the recording medium to protect the image data by these methods, when the digital camera is connected to an external storage device such as a hard disk drive or an information terminal device such as a personal computer including a storage device to store the image data recorded on the recording medium of the digital camera in the external storage device or the storage device arranged in the information terminal device as backup data, the image data is expected to be erased from the recording medium in many cases.\nHowever, in the conventional digital camera, it is very complicated to cancel a read-only attribute, which is set once, because the same troublesome procedure for setting protection described in FIGS. 15A to 15D and FIGS. 16A to 16C must be performed. Therefore, it is also complicated to erase image data, to which a read-only attribute is set, from a recording medium."}
{"text": "Field of the Invention\nThis invention relates to a foldable baby carriage and more particularly to a skeleton construction suitable for large-sized foldable baby carriages.\nAn examination of the history of change of baby carriages with attention paid to their seat sections shows that initially they were of the so-called box type and then changed to the chair type, which has since continued, and baby carriages with chair type seat sections are popular today. The reasons for switching from the box type to the chair type, included a desire for reducing the weight of the baby carriage and for making it more compact in the folded state.\nAlthough the chair type baby carriage is convenient for use, the space allotted to its seat section is limited. As a result, the seat section is reduced in size to the extent that it has been difficult to sit or lay a baby there in a comfortable manner.\nAccordingly, it has been proposed to add a function which makes the baby carriage compact in the closed state and which further increases the area of the seat section when the baby carriage is opened. For example, a function is proposed (for example, in Japanese Patent Publication No. 32067/1963; corresponding to U.S. Pat. No. 4,411,472) which makes the backrest of the seat section reclinable for placing a baby comfortably into the carriage by extending the seat portion of the seat section further forward to allow the baby to stretch his or her legs.\nHowever, even if a backrest reclining mechanism or a seat portion extending function is added, as described above, the space available for the seat section has its limit. This is due mainly to the fact that the body of the baby carriage cannot be made so large because of the design requirement that the baby carriage be made compact when folded. For example, when the length of the seat section extending in the direction of travel of the baby carriage is increased to make it possible for a baby to be comfortably in the carriage lengthwise in the direction of travel of the baby carriage, it is necessary to lengthen the wheel base between the front and rear wheels of the baby carriage in order to stabilize the baby carriage and enhance its safety. Therefore, the members connecting these front and rear wheels are lengthened, resulting in a problem that the closed state of the baby carriage cannot be made so compact."}
{"text": "This invention relates generally to computed tomography (CT) imaging and, more particularly, to methods and apparatus for time-multiplexing data acquisition.\nIn at least one known computed tomography (CT) imaging system configuration, an x-ray source projects a fan-shaped beam which is collimated to lie within an X-Y plane of a Cartesian coordinate system and generally referred to as the xe2x80x9cimaging planexe2x80x9d. The x-ray beam passes through the object being imaged, such as a patient. The beam, after being attenuated by the object, impinges upon an array of radiation detectors. The intensity of the attenuated beam radiation received at the detector array is dependent upon the attenuation of the x-ray beam by the object. Each detector element of the array produces a separate electrical signal that is a measurement of the beam attenuation at the detector location. The attenuation measurements from all the detectors are acquired separately to produce a transmission profile.\nIn known third generation CT systems, the x-ray source and the detector array are rotated with a gantry within the imaging plane and around the object to be imaged so that the angle at which the x-ray beam intersects the object constantly changes. A group of x-ray attenuation measurements, i.e., projection data, from the detector array at one gantry angle is referred to as a xe2x80x9cviewxe2x80x9d. A xe2x80x9cscanxe2x80x9d of the object comprises a set of views made at different gantry angles, or view angles, during one revolution of the x-ray source and detector. In an axial scan, the projection data is processed to construct an image that corresponds to a two dimensional slice taken through the object. One method for reconstructing an image from a set of projection data is referred to in the art as the filtered back projection technique. This process converts the attenuation measurements from a scan into integers called xe2x80x9cCT numbersxe2x80x9d or xe2x80x9cHounsfield unitsxe2x80x9d, which are used to control the brightness of a corresponding pixel on a cathode ray tube display.\nIn a growing number of computer tomography applications it is desirable to increase the speed of data acquisition with corresponding increased patient coverage in a given amount of time. Demands increasingly are put upon the data acquisition system (DAS), particularly in multi-slice applications. To include a single-slice DAS row for each row of detectors is cost-prohibitive. It is known to use a single-slice DAS to process the signal of several combined detector rows. Although this approach may serve to limit the number of DAS systems needed, it negates the potential benefit of higher z-resolution through multiple rows in z.\nIn one exemplary embodiment of the present invention, a single-slice imaging system DAS is time-multiplexed for measuring more than one signal during a DAS standard sampling time interval. More specifically, and in the one embodiment, each detector element includes a scintillator and two photodiodes aligned with respective portions of the scintillator. Each photodiode generates a signal based on the scintillator output. The signal from one of the photodiodes is subjected to a time-dependent gain during consecutive sampling intervals. The gain-adjusted signal and the signal from the other photodiode are combined, and the combined signal is processed to obtain an estimate of the z-derivative of the signal. The estimated z-derivative is then used to generate a high quality image."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device having a circuit in which a thin film transistor (hereinafter referred to as TFT) is formed on a substrate having an insulating surface, and to a manufacturing method thereof. More particularly, the present invention relates to an electro-optical device typically known as a liquid crystal display device provided with an excellent-shaped contact hole, and to electronic equipment with an electro-optical device. Moreover, the semiconductor device according to the present invention refers to all equipment utilizing the semiconductor characteristics for functioning. The above electro-optical device and the electronic equipment with an electro-optical device are also included in semiconductor devices.\n2. Description of the Prior Art\nThe development of a semiconductor device having a large-area integrated circuit formed by the TFT (thin film transistor) on its insulating surface is under progression. AN active matrix liquid crystal display device, an EL display device, a contact-type image sensor, and the like are known as representative examples.\nThe characteristics of the TFT are deteriorated and reliability is lowered when organic resin film is directly formed on the TFT provided on the insulating surface. To solve those problems, conventionally, a laminated organic resin film is formed on the TFT after the formation of an inorganic insulating film (also called passivation film).\nThe TFT is normally connected to wirings through a contact hole. Therefore, when the above inorganic insulating film is provided on the TFT, it is necessary to form a contact hole for connecting to the upper layer wiring in the inorganic insulating film and the organic resin film which covers a TFT gate electrode, a source electrode, or a drain electrode. For instance, the contact hole is formed for connecting a drain electrode of a pixel TFT with a pixel electrode in an active matrix liquid crystal display using TFT.\nA conventional manufacturing process will be described with reference to FIGS. 17A through D. Shown here is an example of an active matrix liquid crystal display applied to this process. Although the pixel TFT is not shown for simplification, a first conductive film 11 is identical with a drain electrode of a pixel TFT or electrically connected thereto. Also not shown is that there is a single layer or a multiple layer of insulating film between a substrate 10 and the first conductive film 11.\nShown in FIG. 17A is a state where the first conductive film 11 identical with the drain electrode of a pixel TFT or electrically connected thereto is formed on the substrate 10 on which the pixel TFT is formed. An inorganic insulating film 12 and an organic resin film 13 formed in laminations are formed on top of the first conductive film 11.\nAnd shown in FIG. 17B is a state where the first patterning is performed by using a resist mask 14 by way of photolithography, and a first contact hole is opened only in the organic resin film 13.\nSubsequently, a second patterning is performed by using a resist mask 15 after removing the resist mask 14, and a second contact hole is opened only in the inorganic insulating film 12. This state is shown in FIG. 17C. Since the second contact hole is formed in the bottom portion of the first contact hole, the diameter of its opening is smaller than that of the first contact hole.\nShown in FIG. 17D next is a state where a pixel electrode 16, made of transparent conductive film, is formed after removing the resist mask 15.\nAs shown in FIG. 17D, in this way there is a step in the shape of the contact hole because it has been formed after the first and second patterning processes.\nAdditionally, besides the above conventional manufacturing method, another method is to perform patterning right after forming the inorganic insulating film, then form the organic resin film and perform patterning again to form a contact hole. Two patterning processes were also necessary even in this method.\nSince the number of processes and masks has increased due to two patterning (organic resin film patterning and inorganic insulating film patterning) processes in the conventional method, this led to an increase in costs.\nIn the two patterning processes, each method uses different photo mask, and therefore poor contact occurred when the masks had not been overlapped in precision. Also, in the example of the conventional process shown in FIG. 17, fining of the contact hole is difficult. That is because considering the margin when overlapping, the opening diameter of the second contact hole that was opened in the second patterning is 1.5 to 2 times bigger than the opening diameter of the first contact hole opened in the first patterning process.\nFurthermore, the shape of a conventional contact hole (of which an example is shown in FIG. 17D) is a complicated shape formed by overlapping two contact holes with different opening diameter. Thus, poor coverage has occurred on a second conductive film formed later.\nA technique of the present invention is for solving the above problems, and therefore it is an object of the present invention to provide a manufacturing method of a semiconductor device whereby the number of processes is decreased due to simultaneously forming a contact hole in a lamination film (inorganic insulating film and organic resin film) of different material and film thickness by conducting etching once.\nIt is another object of the present invention to improve operating efficiency and reliability of a semiconductor device by providing a contact hole that is uniform in shape, and moreover an appropriate one.\nIt is still another further object of the present invention to form a pixel electrode of good coverage and to provide a structure for improving the yield of an active matrix type liquid crystal display device.\nIn order to solve the above problems, the present invention provides a semiconductor device comprising: a first conductive film formed on an insulating substrate; an inorganic insulating film covering said first conductive film; an organic resin film covering said inorganic insulating film; a contact hole that goes through said inorganic insulating film and said organic resin film; and a second conductive film formed on said organic resin film which is connected to said first conductive film at a bottom surface of said contact hole.\nFurther, according to the above structure, said contact hole is formed by performing one etching.\nStill further, according to each structure of the above, an edge portion of an inorganic insulating film that comes in contact with a bottom surface of said contact hole is taper like having an angle range of 30xc2x0 to 800xc2x0 from a horizontal surface.\nFurther, according to each structure of the above, an edge portion of an organic resin film that comes in contact with said inorganic insulating film has an angle range of 50xc2x0 to 90xc2x0 from a horizontal surface.\nStill further, according to each structure of the above, a TFT is electrically connected to said first conductive film.\nFurther, according to each structure of the above, said second conductive film is a pixel electrode.\nAccording to each structure of the above, said inorganic insulating film is a silicon nitride film or a silicon oxide nitride film.\nMoreover, in order to realize the above structure, the present invention provides a method of manufacturing a semiconductor device comprising the steps of: forming a first conductive film; forming an inorganic insulating film on said first conductive film; forming an organic resin film on said inorganic insulating film; forming a contact hole in a laminated film formed of said inorganic insulating film and said organic resin film in one process; and forming a second conductive film in said contact hole.\nFurther, according to the above structure, said process of forming a contact hole is performed by dry etching employing mixed gas containing fluorine-based etchant gas and oxygen gas.\nFurthermore, according to the above structure, a selective ratio of an etching rate of said inorganic insulating film to an etching rate of said organic resin film is 1.6 to 2.9.\nStill further, according to the above structure, said inorganic insulating film is a silicon nitride film or a silicon oxide nitride film."}
{"text": "In eukaryotic cells triacylglycerols are quantitatively the most important storage form of energy. The main pathway for synthesis of triacylglycerol is believed to involve three sequential acyl-transfers from acyl-CoA to the glycerol backbone. Acyl-CoA:diacylglycerol acyltransferase (DAGAT, EC 2.3.1.20) uses fatty acyl-CoA (acyl donor) and 1,2-diacylglycerol (acyl acceptor) as substrates to catalyze the third and only committed step in triacylglycerol synthesis. DAGAT plays a fundamental role in the metabolism of cellular glycerolipids.\nUntil recently, it was believed that only DAGAT could carry out the final step in triacylglycerol biosynthesis. However, it has now been demonstrated that microsomal preparation of developing seeds from several plants (sunflower, Helianthus annus; castor bean, Ricinus communis; hawk's beard, Crepis palaestina) as well as microsomal preparations from yeast (Saccharomyces cerevisiae) catalyze triacylglycerol formation via the enzyme phospholipid:diacylglycerol acyltransferase (PDAT, EC 2.3.1.158, Registry Number 288587-47-3) (WO 2000060095; Dahlqvist et al., Proc. Natl. Acad. Sci. USA 97(12):6487-6492 (2000)). This enzyme differs from DAGAT by synthesising triacylglycerol using an acyl-CoA-independent mechanism. The specificity of the enzyme for the acyl group in the phospholipid varies with species, e.g., the enzyme from castor bean preferentially incorporates vernoloyl (12,13-epoxyoctadec-9-enoyl) groups into triacylglycerol, whereas that from the hawk's beard incorporates both ricinoleoyl (12-hydroxyoctadec-9-enoyl) and vernoloyl groups. The enzyme from the yeast Saccharomyces cerevisiae specifically transfers acyl groups from the sn-2 position of the phospholipid to diacylglycerol, thus forming an sn-1-lysophospholipid. It has also been shown that PDAT activity is present in vegetative tissues of Arabidopsis thaliana (Banaś et al., Biochem. Soc. Trans. 28:703-703 (2000)). Furthermore, the substrate specificity of PDAT varies between species and depends on the head group of the acyl donor, the acyl group transferred and the acyl chains of the acyl acceptor (1,2-diacylglycerol)."}
{"text": "The invention relates to a clamping system for crankshafts and similar parts, on crankshaft working machines.\nIn clamping-up crankshafts and similar workpieces in machines for working them, especially when they have been forged or cast and have not yet been machined, the problem consists particularly in mounting them firmly in such a definite and reproducible manner that, in spite of the relatively small cross section of a crankshaft in comparison to its total length, no flexing or twisting of the crankshaft will occur in spite of the cutting forces and temperatures involved in the machining.\nThe manufacturing process most frequently employed at the present time is to work the crankshaft casting, before it is actually machined on a crankshaft cutting machine, at one or more cheek surfaces precisely defined with reference to the crankshaft axis of rotation, using a special machine for the purpose. In the subsequent process on the crankshaft cutting machine, the crankshaft is then clamped at the previously worked cheek surfaces by means of levers and faceplates, so that any inaccuracy of these surfaces with respect to the axis of rotation of the crankshaft results in a distortion of the crankshaft when it is clamped in the crankshaft working machine. This, together with the deformation of the crankshaft due to the cutting forces, expresses itself in an off-center position of the machined journal areas with reference to the crankshaft axis after it is unclamped. The disadvantage of this manner of proceeding consists mainly in the fact that the working of a crankshaft in this manner is very complicated and time-consuming, since on the one hand a special machine is needed for the production of the pre-cut clamping surfaces. Furthermore, it is found to be disadvantageous that, in the case of another necessary clamping-up for another series of operations, such as a finishing operation or clearance cutting operation, the crankshaft is put under tension because the crankshaft is deformed or distorted as a result of the preceding roughing operation, and thus the pre-cut clamping surfaces are no longer in the originally defined position with respect to the axis of rotation of the crankshaft. This error then has a negative effect accordingly on the subsequent machining accuracy. This necessitates greater machining allowances and hence increased cost in the next manufacturing operation, such as the grinding operation for example. Furthermore, the tension control that is essential in a clamping device of this kind, is not available.\nThe invention is therefore addressed to the task of creating a clamping device by which crankshafts which have not yet been machined on the journals and cheeks, can be clamped up in crankshaft processing machines in a defined position repeatably and without flexing the crankshaft, even in the case of slightly differing crankshaft forgings and varying clamping forces."}
{"text": "A typical production-packaging apparatus for wrapping a stack of objects with a film has at least one film supply for the film, a film-feed head for feeding the film, and a device for pulling the film down over the stack of objects. The invention further relates to a method of wrapping a stack of objects with a film. Within the scope of the invention, “film” means in particular a plastic film, preferably an elastic thermoplastic film. It is within the scope of the invention that the stack of objects to be wrapped is a stack of products. The stack may also be a single object or whiteware, for example refrigerators or the like. The stack of objects to be wrapped is advantageously accommodated on a pallet.\nIn practice, it is not uncommon for stacks of objects or pallets carrying these stacks of objects having different orientations to be consecutively supplied to a wrapping station of an apparatus for wrapping the stack of objects. Thus, a stack of objects that in top view are not square or a pallet that in top view is not square having the length/width dimensions L/W is initially supplied and wrapped, and subsequently another stack of objects that has the dimensions W/L and that is rotated by 90° about a vertical axis with respect to the first stack of objects must be supplied and wrapped. Additional measures are necessary here in order to obtain a film or film sheet that in such a case is also optimally adapted to the dimensions of the stack of objects or of the pallet.\nIn practice it is known, among other things, to lift the stack of objects using the elevator of a shrink-wrapping device and to rotate the stack of objects by 90° for the covering operation. However, this often results the stack of objects falling over, together with product losses and facility down times. It is also known to provide in the conveyor a rotatable table underneath the film-feed head so that the stack of objects may be appropriately rotated. However, constrictions result in the transition regions between the rotatable table and the conveyor edge due to the fact that the pallets associated with the stacks of objects often become stuck at the transitions, and the stacks of objects then frequently fall over. It is also known in practice to design the entire film-feed head to be rotatable by 90° so that the film can always be fed in the desired direction. This approach is relatively complicated, and the process is also too slow, in particular because of the large mass that has to be moved.\nIt is further known in practice for the film initially to be completely provided in the feed head and welded, and then transferred to a roller and rolled up by the roller. This roller is then rotated by 90° with respect to the direction of transport, and the film is transferred from the roller to the pull-down device. These measures known in practice are also relatively complicated. The multiple transfers of the film result in a relatively high susceptibility to malfunctions."}
{"text": "Many different solvents, including water and organic solvents are used in coating, printing, dyeing, impregnating, sizing, papermaking, waterproofing, waxing or otherwise treating webs made of paper, plastic, metal, cloth, rubber, textiles, etc. The principal object of the invention is to provide an improved method and apparatus for drying such webs. It is a further object to provide a method and apparatus for recovering the solvent for re-use.\nIn the past, forced circulation of air over the surface of a web within closed vacuum chambers has been employed to dry the web as disclosed in U.S. Pat. No. 1,487,362 to O. D. Rice. It has also been known to dry the web by heating it with radiant heaters or gas flames which when organic solvents are used can also incinerate the organic solvent vapor. In the past, the recovery of a portion of the organic solvent has been accomplished by passing the solvent vapors removed from the web through a series of absorbent activated carbon beds, then passing steam through the solvent saturated activated carbon bed to remove the absorbed solvent as solvent-water mixture, which was subsequently in a secondary process fractionally distilled to separate the organic solvent from the water. However, the activated carbon beds are limited to absorbing only a small percentage of their weight in organic solvent, and, therefore, the beds must be quite large. Several of these large beds are typically required, inasmuch as the organic solvent vapors must be passed through an unsaturated bed while the organic solvent is being removed from the saturated bed or beds. In addition to large size, the activated carbon absorption method is relatively expensive in cost and use of energy and recovers only a portion of the solvent used in processing the web."}
{"text": "1. Field of the Invention\nThe present invention relates to an image capturing apparatus, a control method of the image capturing apparatus, and an image capturing system, and more specifically, to an image capturing apparatus having auto focus control function, a control method of the image capturing apparatus, and an image capturing system.\n2. Description of the Related Art\nTypical schemes of focus control methods for image capturing apparatuses include a contrast AF scheme and a phase difference AF scheme. The contrast AF scheme and the phase difference AF scheme are AF schemes widely used in video cameras and digital still cameras, with an image sensor being used as a focus detection sensor. In these focus control methods, focus detection results may contain errors due to various aberrations of an optical system. Various methods have been proposed to reduce such errors.\nFor example, Japanese Patent No. 5087077 discloses a method of calculating a correction value for correcting a focus detection result in accordance with a frequency band for evaluation of a signal used for focus detection. Regardless of the above-mentioned focus control methods, such a focus detection error occurs depending on, for example, a frequency band for evaluation of a focus control signal used in the contrast AF scheme or the phase difference AF scheme.\nWith the structure described in Japanese Patent No. 5087077, however, when the number of kinds of signals used for focus detection is larger, more correction values are needed. This requires more memory to store the correction values, leading to an increase in cost. Especially in the case of an image capturing apparatus having auto focus control function to which an interchangeable imaging lens unit can be attached, the imaging lens unit needs to store correction values corresponding to different focus detection functions or different focus detection methods, and the need for more memory leads to an increase in cost of the imaging lens unit."}
{"text": "The semiconductor integrated circuit (IC) industry has experienced rapid growth. Technological advances in IC materials and design have produced generations of ICs. Each generation has smaller and more complex circuits than the previous generation.\nIn the course of IC evolution, functional density (i.e., the number of interconnected devices per chip area) has generally increased while geometric size (i.e., the smallest component (or line) that can be created using a fabrication process) has decreased. This scaling-down process generally provides benefits by increasing production efficiency and lowering associated costs.\nSuch scaling down has also increased the complexity of processing and manufacturing ICs and, for these advances to be realized, similar developments in IC processing and manufacturing are needed. For example, a three-dimensional transistor, such as a semiconductor device with fin field-effect transistors (FinFETs), has been introduced to replace planar transistors. These relatively new types of semiconductor IC devices face manufacturing challenges, and they have not been entirely satisfactory in all respects."}
{"text": "Identification of underground hydrocarbon bearing formations, and extraction of hydrocarbons from such formations, generates a significant amount of data. In particular, each borehole drilled into a hydrocarbon bearing formation may be logged by a host of different logging tools, both during drilling, and after the borehole is cased. Generally speaking, the data and information gathered about the earth formations surrounding the borehole are stored, and in many cases the data and information are used in the planning and drilling of other boreholes in relatively close proximity.\nConsider a situation where multiple boreholes have been previously drilled into a hydrocarbon bearing formation, and multiple logs have been taken within each borehole. When planning the next borehole to be drilled, the formation properties along the proposed borehole path will be estimated using the logs from all the multiple previously drilled boreholes. The data resolution for each log may on the order of six inches (i.e., a datum in the log representing the value of the formation parameter measured for every six inches), and thus the number of data points to consider when estimating the formation properties along the proposed borehole path is enormous. For this reason, in the related art updating of models or predictions of formation parameters along a proposed borehole path cannot be accomplished in real time with the drilling along the proposed borehole path.\nThus, any advance which results in an ability to calculate updates of predicted formation parameters in less time or with less computing power, would provide a competitive advantage."}
{"text": "1. Field of the Invention\nThe present invention relates to a personal watercraft.\n2. Description of the Related Art\nA conventional personal watercraft is described in Japanese Unexamined Utility Model Application Publication No. H03-088997 (1991). This personal watercraft includes a storage portion provided at the stern, a lid arranged to open and close the storage portion, and a rope to be used when a passenger gets on and off the personal watercraft. The rope is stored in the storage portion.\nA conventional personal watercraft is described in U.S. Patent Application Publication No. 2009/0158982 A1. This personal watercraft includes a sub deck and a platform disposed at the stern. The rear portion of the sub deck defines a storage portion. The platform is disposed on the sub deck. The platform is joined to the sub deck so as to be rotatable up and down. The storage portion is opened by the platform."}
{"text": "When a great amount of the raw water containing the SS at a high concentration as mentioned above is to be filtered, an amount of the SS to be captured in a filtering layer needs to be increased. For this purpose, the filtration is commonly carried out in a down-flow system by using a precipitatable granular filter medium having particle diameters of around 10 mm, which is greater than those employed in an ordinary sand filtration. However, if the diameters of the filter medium are large, the SS-removing percentage inevitably decreases. Further, since the SS is captured by the entire filtering layer by back washing, there are problems in that it is difficult to discharge the SS captured, and a great amount of wash water is required in the back washing.\nWhen an up-flow filtration system is employed, filter media made of polypropylene or polyethylene having a density of 0.9 to 0.95 g/cm3 each in the form of cylindrical pieces is often used (In this application, the term “medium” means a single small piece, and the term “media” means a mass of a number of such small pieces or mediums). However, since such filter media have a small difference in density between water, there is a fear that the filter media having cylindrical forms flow out. In view of this, a screen is provided under a filter media-filled layer. In this case, the screen is clogged after being repeatedly subjected to the filtration and back washing operations, which disenables the filtration. Furthermore, it is very difficult from the standpoint of the structure of the apparatus to eliminate the clogging of the lower screen.\nFurther, since the screen exists during back washing, the contaminant clogging the screen cannot be discharged and removed through gaps being formed among the filtering media due to their expansion. Therefore, the filter media need to be repeatedly washed with air and water, which requires a long-time back washing step. In addition, the mesh of the screen is often generally increased to as much as 10 mm or so in order to prevent the lower screen from being clogged, the grain diameters of the filter media need to be greater than that mesh, and the thickness of the filter media needs to be larger than 2 m or so in order to enhance the SS-capturing efficiency. In this case, the apparatus unfavorably becomes very large."}
{"text": "This specification relates to a data communication system and, in particular, a system that implements real-time, scalable publish-subscribe messaging.\nThe publish-subscribe pattern (or “PubSub”) is a data communication messaging arrangement implemented by software systems where so-called publishers publish messages to topics and so-called subscribers receive the messages pertaining to particular topics to which they are subscribed. There can be one or more publishers per topic and publishers generally have no knowledge of what subscribers, if any, will receive the published messages. Some PubSub systems do not cache messages or have small caches meaning that subscribers may not receive messages that were published before the time of subscription to a particular topic. PubSub systems can be susceptible to performance instability during surges of message publications or as the number of subscribers to a particular topic increases."}
{"text": "(a) Field of the Invention\nThe invention relates to an automatic noninvasive blood pressure monitor, i.e. an instrument for measuring the blood pressure in man on arteries being compressible from the surface by means of a pressure cuff or a pelotte equipped with a plethysmographic gauge, such as an impedance or a photo-electric one, which gauge is connected through at least one amplifier and a phase corrector to an electro-pressure transducer. All these components constitute the closed loop of a servosystem which continuously and instantaneously changes the pressure in the cuff and maintains thus the volume of the artery at a value corresponding to the zero tension of the arterial wall. The pressure within the cuff thus follows the instantaneous value of intraarterial pressure. In such instruments automatic setting and correction of the setpoint and gain is also possible.\n(b) Description of the Prior Art\nSimilar instruments are known from the patent literature, e.g. Czechoslovak patent specification No. 133205, U.S. patent specification No. 4,510,940, as well as from articles in medical and technical journals. The instruments known up to now, however, either have no automatic initial setting of the setpoint and gain, or have no correction of these parameters during the measurement. Although an instrument with such a correction has been proposed, the correction is performed during short-lasting interruption of the measurement. Another instrument proposed is equipped with an additional cuff placed on another circulatory region so that a pair of cuffs is necessary. The proposed instruments permit measurement of the blood pressure on arteries of the finger or of another region which can easily be transilluminated, however continuous blood pressure measuring is impossible on other arteries especially on larger ones."}
{"text": "1. Field of the Invention\nThis invention concerns a semiconductor device used for mobile communication apparatus such as cellulars operation in a microwave band of 500 MHz to 2.5 GHz and, more in particular, it relates to a technique which is effective to high frequency powers amplifier of power-amplifying high frequency signals and outputting them.\n2. Description of Related Art\nIn recent years, mobile communication apparatus typically represented by communication systems such as GSM (Global System for Mobile Communication), PCS (Personal Communication system), PDC (Personal Digital Cellular) system and CDMA (Code Division Multiple Access) system (so-called hand set telephone) have been popularized world wide.\nGenerally, a mobile communication apparatus comprises an antenna for emitting and receiving electromagnetic waves, a high frequency power amplifier for amplifying high frequency signals under power modulation and supplying them to the antenna, a receiving section for processing high frequency signals received on the antenna, a control section for controlling them and a battery for supplying a power source voltage to them.\nSuch mobile communication apparatus and semiconductor devices used for the mobile communication apparatus have been disclosed, for example, in the known literatures described below.\n(1) The constitution of the mobile communication apparatus is disclosed, for example, in “HITACHI HYORON” vol. 78, No. 11 (1996-11), pages 21 to 26 (Literature 1).\n(2) Constitution of a typical high frequency power amplifier on the GSM system is described, for example, in ISSCC 98, DIGEST OF TECHNICAL PAPERS (Feb. 5, 1998), pages 50 to 55 (Literature 2).\nAccording to the literatures, a threshold voltage of FET is controlled to an appropriate level for stable circuit design and decrease of leak current in an off state. In the constitution of the amplifier, two chips are arranged in parallel for final stage devices in a 3-stage amplifier circuit and synthesizing and an impedance-matching circuit is disposed to each of them, for the synthesis of outputs of attain higher output than in the case of using one chip. The constitution of this amplifier is referred to as DD-DIMA technique (Divided Device and Collectively Impedance Matched Amplifier) in the literature.\n(3) Amplifying devices applied to the high frequency power amplifiers are described, for example, in IEDM 97 Technical Digest (1997), pages 51 to 54 (Literature 3).\nThis literature discloses that an amplifier device is constituted with a power insulated gate type field effect transistor (hereinafter simply referred to as power MOSFET) using Si (silicon) semiconductors to attain high performance.\nSpecifically, the performance is improved by defining the gate length of a MOSFET to 0.4 μm. Further, the drain breakdown voltage is set to 20 V or higher by disposing an offset layer of an about 0.7 μm length on the side of the drain of the power MOSFET. Further, it is important to lower the gate resistance for high frequency operation and the gate resistance is decreased by a structure of short circuiting a metal silicide/silicon laminated gate electrode with an aluminum wiring (Al-shorted silicon gate structure).\n(4) There has been a trend of adopting compound semiconductor (GaAs) wafers for making the device efficiency higher. Such a technical trend is described, for example, in NIKKEI ELECTRONICS 1998, 11, 2 (No. 729), pages 238 to 245 (Literature 4). However, as described also in this literature, the unit wafer cost in the GaAs technique is more expensive compared with Si.\nFor generalized use of the mobile communication apparatus, it has been demanded for further reduction of the size and weight, and decrease in the power consumption of the apparatus. For this purpose, it is necessary to attain further reduction in the size and weight and decrease in the power consumption for each of components constituting the mobile communication apparatus.\nOne of the components is a high frequency power amplifier for supplying high frequency signals to the antenna. Generally, the high frequency power amplifier most consumes electric power and it is effective of saving the consumption power of the mobile apparatus to decrease the consumption power of the high frequency power amplifier (improvement of efficiency).\nIn the GSM system amplifier using silicon (Si) semiconductors, it was attained an output voltage of about 3.5 W and a overall efficiency (η all) of about 50% at a working frequency of 900 MHz and a power source voltage of 3.5 V. The overall efficiency is an efficiency in a high frequency power amplifier constituted with a three stage power amplifier of power MOSFET (high frequency module).\nThe performance of the power MOSFET using Si as the output stage amplifier device is about 55% power of a power-added efficiency (η add) at 2 W output based on the DD-CIMA technique, and it was necessary to attain a power added efficiency of 65% or more in power MOSFET in order to improve the overall efficiency of the amplifier to 55% or more.\nThe power added efficiency (η add) in the microwave power MOSFET is defined, for example, in “Optical Microwave Semiconductor Applied Technology” Feb. 29, 1996, first edition, first print (Published from Science Forum Co.), pages 59 to 66 (Literature 5).\nAlso in the PCS system amplifier, an output voltage of 2 W and an overall efficiency of 45% at a working frequency of 1900 MHz have been attained. The performance of the power MOSFET as the output stage amplifier device is about 50% at 1 W output. For improving the overall efficiency of the amplifier to 50% or more, it was necessary to attain the power-added efficiency of 55% or higher in the power MOSFET.\nFor improving the power-added efficiency of the amplifier device (power MOSFET), it is considered to decrease the on-resistance, gate resistance and parasitic capacitance and improvement of the mutual conductance."}
{"text": "A radio frequency (RF) input signal is typically applied to an input of a home entertainment device, for example, a set-top box via an external coaxial cable, for example, a Multimedia over Coax Alliance (MoCA) cable, having a terminating coaxial connector known as F connector. The set-top box contains electronic components which may be vulnerable to damage from electrical surges. Surges in transmission line voltages, as well as noise, can change the operating range of the components and severely damage and/or destroy them. There are many sources which can cause harmful electrical energy surges. One source is RF interference that can be coupled to transmission lines from a multitude of sources. The transmission lines act as large antennas that may extend over several miles, thereby collecting a significant amount of RF noise power from such sources as radio broadcast antennas. Another source could be lightening. Therefore, it may be desirable to interpose an RF surge suppression device, forming direct current (DC) short or low impedance when a surge occurs. The RF surge suppression device may be coupled to an output of the F connector.\nIt may be desirable to have the F connector mounted at some distance above a printed circuit board (PCB) of the set-top box, for example, to facilitate easier manufacturing and/or more convenient placement or usage of a lower cost implementation of the connector. In such arrangement, the lead that goes from a center conductor of the connector to a pad in the PCB, the rest of the F connector and the RF surge suppression device may exhibit a reactance that may not be negligible for impedance matching purposes. Additionally, at the PCB, the pad may be relatively large to facilitate machine or reflow soldering. This pad may introduce a stray capacitance to ground that also may not be negligible.\nA receiver, such as the set-top box, may receive through the connector an input signal transmitted in a multiple transmission mode system covering a wide range of frequencies, for example, up to 2.15 GHz. It may be desirable to avoid significant input return loss as a result of the aforementioned reactance of the network that includes the connector and the RF surge suppression device. Return loss is a measurement of how well the impedance of a load, including, for example, a filter that is driven by the network, including the connector and the RF surge suppression device, is matched to the characteristic impedance of the network. The return loss is a number associated with a corresponding interface that is calculated from the reflection caused at the corresponding interface as a result of an impedance mismatch. The return loss is usually expressed as a ratio in decibels (dB)."}
{"text": "1. Field of Invention\nEmbodiments of the invention relate generally to pattern-recognition processors and, more specifically, in certain embodiments, to pattern-recognition processors having local input conductors.\n2. Description of Related Art\nIn the field of computing, pattern recognition tasks are increasingly challenging. Ever larger volumes of data are transmitted between computers, and the number of patterns that users wish to identify is increasing. For example, spam or malware are often detected by searching for patterns in a data stream, e.g., particular phrases or pieces of code. The number of patterns increases with the variety of spam and malware, as new patterns may be implemented to search for new variants. Searching a data stream for each of these patterns can form a computing bottleneck. Often, as the data stream is received, it is searched for each pattern, one at a time. The delay before the system is ready to search the next portion of the data stream increases with the number of patterns. Thus, pattern recognition may slow the receipt of data.\nHardware has been designed to search a data stream for patterns, but this hardware often consumes excessive amounts of power. Some devices configured to search a data stream do so by distributing the data stream among a plurality of circuits. The circuits each determine whether the data stream matches a portion of a pattern. Often, a large number of circuits operate in parallel, each searching the data stream at generally the same time. Each of these circuits consumes electrical energy and generates heat. As a result, the amount of heat generated and the amount of electricity consumed scales with the number of patterns. In devices with a large number of circuits, the cost of the electricity and the difficulty of removing this heat complicate the use of pattern-recognition hardware."}
{"text": "This disclosure relates generally to tracking devices, and more specifically, to modifying the behavior of a tracking device or an associated mobile device based on its presence in a zone or an area.\nElectronic tracking devices have created numerous ways for people to track the locations of people and/or objects. For example, a user can use GPS technology to track a device remotely or determine a location of the user. In another example, a user can attach a tracking device to an important object, such as keys or a wallet, and use the features of the tracking device to more quickly locate the object, (e.g., if it becomes lost).\nThe figures depict various embodiments of the present invention for purposes of illustration only. One skilled in the art will readily recognize from the following discussion that alternative embodiments of the structures and methods illustrated herein may be employed without departing from the principles of the invention described herein."}
{"text": "Mobile communications has changed the way people communicate and mobile phones have been transformed from a luxury item to an essential part of every day life. The use of mobile phones is today dictated by social situations, rather than hampered by location or technology. While voice connections fulfill the basic need to communicate, and mobile voice connections continue to filter even further into the fabric of every day life, the mobile Internet is the next step in the mobile communication revolution. The mobile Internet is poised to become a common source of everyday information, and easy, versatile mobile access to this data will be taken for granted.\nThird generation (3G) cellular networks have been specifically designed to fulfill these future demands of the mobile Internet. As these services grow in popularity and usage, factors such as cost efficient optimization of network capacity and quality of service (QoS) will become even more essential to cellular operators than it is today. These factors may be achieved with careful network planning and operation, improvements in transmission methods, and advances in receiver techniques. To this end, carriers need technologies that will allow them to increase downlink throughput and, in turn, offer advanced QoS capabilities and speeds that rival those delivered by cable modem and/or DSL service providers. In this regard, networks based on wideband CDMA (WCDMA) technology may make the delivery of data to end users a more feasible option for today's wireless carriers.\nFIG. 1A is a technology timeline indicating evolution of existing WCDMA specification to provide increased downlink throughput. Referring to FIG. 1A, there is shown data rate spaces occupied by various wireless technologies, including General Packet Radio Service (GPRS) 100, Enhanced Data rates for GSM (Global System for Mobile communications) Evolution (EDGE) 102, Universal Mobile Telecommunications System (UMTS) 104, and High Speed Downlink Packet Access (HSDPA) 106.\nThe GPRS and EDGE technologies may be utilized for enhancing the data throughput of present second generation (2G) systems such as GSM. The GSM technology may support data rates of up to 14.4 kilobits per second (Kbps), while the GPRS technology, introduced in 2001, may support data rates of up to 115 Kbps by allowing up to 8 data time slots per time division multiple access (TDMA) frame. The GSM technology, by contrast, may allow one data time slot per TDMA frame. The EDGE technology, introduced in 2003, may support data rates of up to 384 Kbps. The EDGE technology may utilizes 8 phase shift keying (8-PSK) modulation for providing higher data rates than those that may be achieved by GPRS technology. The GPRS and EDGE technologies may be referred to as “2.5G” technologies.\nThe UMTS technology, introduced in 2003, with theoretical data rates as high as 2 Mbps, is an adaptation of the WCDMA 3G system by GSM. One reason for the high data rates that may be achieved by UMTS technology stems from the 5 MHz WCDMA channel bandwidths versus the 200 KHz GSM channel bandwidths. The HSDPA technology is an Internet protocol (IP) based service, oriented for data communications, which adapts WCDMA to support data transfer rates on the order of 10 megabits per second (Mbits/s). Developed by the 3G Partnership Project (3GPP) group, the HSDPA technology achieves higher data rates through a plurality of methods. For example, many transmission decisions may be made at the base station level, which is much closer to the user equipment as opposed to being made at a mobile switching center or office. These may include decisions about the scheduling of data to be transmitted, when data is to be retransmitted, and assessments about the quality of the transmission channel. The HSDPA technology may also utilize variable coding rates. The HSDPA technology may also support 16-level quadrature amplitude modulation (16-QAM) over a high-speed downlink shared channel (HS-DSCH), which permits a plurality of users to share an air interface channel\nIn some instances, HSDPA may provide a two-fold improvement in network capacity as well as data speeds up to five times (over 10 Mbit/s) higher than those in even the most advanced 3G networks. HSDPA may also shorten the roundtrip time between network and terminal, while reducing variances in downlink transmission delay. These performance advances may translate directly into improved network performance and higher subscriber satisfaction. Since HSDPA is an extension of the GSM family, it also builds directly on the economies of scale offered by the world's most popular mobile technology. HSDPA may offer breakthrough advances in WCDMA network packet data capacity, enhanced spectral and radio access networks (RAN) hardware efficiencies, and streamlined network implementations. Those improvements may directly translate into lower cost-per-bit, faster and more available services, and a network that is positioned to compete more effectively in the data-centric markets of the future.\nThe capacity, quality and cost/performance advantages of HSDPA yield measurable benefits for network operators, and, in turn, their subscribers. For operators, this backwards-compatible upgrade to current WCDMA networks is a logical and cost-efficient next step in network evolution. When deployed, HSDPA may co-exist on the same carrier as the current WCDMA Release 99 services, allowing operators to introduce greater capacity and higher data speeds into existing WCDMA networks. Operators may leverage this solution to support a considerably higher number of high data rate users on a single radio carrier. HSDPA makes true mass-market mobile IP multimedia possible and will drive the consumption of data-heavy services while at the same time reducing the cost-per-bit of service delivery, thus boosting both revenue and bottom-line network profits. For data-hungry mobile subscribers, the performance advantages of HSDPA may translate into shorter service response times, less delay and faster perceived connections. Users may also download packet-data over HSDPA while conducting a simultaneous speech call.\nHSDPA may provide a number of significant performance improvements when compared to previous or alternative technologies. For example, HSDPA extends the WCDMA bit rates up to 10 Mbps, achieving higher theoretical peak rates with higher-order modulation (16-QAM) and with adaptive coding and modulation schemes. The maximum QPSK bit rate is 5.3 Mbit/s and 10.7 Mbits with 16-QAM. Theoretical bit rates of up to 14.4 Mbit/s may be achieved with no channel coding. The terminal capability classes range from 900 kbit/s to 1.8 Mbit/s with QPSK modulation, and 3.6 Mbit/s and up with 16-QAM modulation. The highest capability class supports the maximum theoretical bit rate of 14.4 Mbit/s.\nHowever, implementing advanced wireless technologies such as WCDMA and/or HSDPA may still require overcoming some architectural hurdles. For example, the RAKE receiver is the most commonly used receiver in CDMA systems, mainly due to its simplicity and reasonable performance and WCDMA Release 99 networks are designed so that RAKE receivers may be used. A RAKE receiver contains a bank of spreading sequence correlators, each receiving an individual multipath. A RAKE receiver operates on multiple discrete paths. The received multipath signals may be combined in several ways, from which maximum ratio combining (MRC) is preferred in a coherent receiver. However, a RAKE receiver may be suboptimal in many practical systems, for example, its performance may degrade from multiple access interference (MAI), that is, interference induced by other users in the network.\nIn the case of a WCDMA downlink, MAI may result from inter-cell and intracell interference. The signals from neighboring base stations compose intercell interference, which is characterized by scrambling codes, channels and angles of arrivals different from the desired base station signal. Spatial equalization may be utilized to suppress inter-cell interference. In a synchronous downlink application, employing orthogonal spreading codes, intra-cell interference may be caused by multipath propagation. Due to the non-zero cross-correlation between spreading sequences with arbitrary time shifts, there is interference between propagation paths (or RAKE fingers) after despreading, causing MAI. The level of intra-cell interference depends strongly on the channel response. In nearly flat fading channels, the physical channels remain almost completely orthogonal and intra-cell interference does not have any significant impact on the receiver performance. On the other hand, the performance of the RAKE receiver may be severely deteriorated by intra-cell interference in frequency selective channels. Frequency selectivity is common for the channels in WCDMA networks.\nDue to the difficulties faced when non-linear channel equalizers are applied to the WCDMA downlink, detection of the desired physical channel with a non-linear equalizer may result in implementing an interference canceller or optimal multi-user receiver. Both types of receivers may be prohibitively complex for mobile terminals and may require information not readily available at the mobile terminal. Alternatively, the total base station signal may be considered as the desired signal. However, non-linear equalizers rely on prior knowledge of the constellation of the desired signal, and this information is not readily available at the WCDMA terminal. The constellation of the total base station signal, that is, sum of all physical channels, is a high order quadrature amplitude modulation (QAM) constellation with uneven spacing. The spacing of the constellation changes constantly due to transmission power control (TPC) and possible power offsets between the control data fields, time-multiplexed to the dedicated physical channels. The constellation order may also frequently change due to discontinuous transmission. This makes an accurate estimation of the constellation very difficult.\nIn this regard, the use of multiple transmit and/or receive antennas may result in an improved overall system performance. These multi-antenna configurations, also known as smart antenna techniques, may be utilized to mitigate the negative effects of multipath and/or signal interference on signal reception. It is anticipated that smart antenna techniques may be increasingly utilized both in connection with the deployment of base station infrastructure and mobile subscriber units in cellular systems to address the increasing capacity demands being placed on those systems. These demands arise, in part, from a shift underway from current voice-based services to next-generation wireless multimedia services that provide voice, video, and data communication.\nThe utilization of multiple transmit and/or receive antennas is designed to introduce a diversity gain and to suppress interference generated within the signal reception process. Such diversity gains improve system performance by increasing received signal-to-noise ratio, by providing more robustness against signal interference, and/or by permitting greater frequency reuse for higher capacity. In communication systems that incorporate multi-antenna receivers, a set of M receive antennas may be utilized to null the effect of (M−1) interferers, for example. Accordingly, N signals may be simultaneously transmitted in the same bandwidth using N transmit antennas, with the transmitted signal then being separated into N respective signals by way of a set of N antennas deployed at the receiver. Systems that utilize multiple transmit and receive antennas may be referred to as multiple-input multiple-output (MIMO) systems. One attractive aspect of multi-antenna systems, in particular MIMO systems, is the significant increase in system capacity that may be achieved by utilizing these transmission configurations. For a fixed overall transmitted power, the capacity offered by a MIMO configuration may scale with the increased signal-to-noise ratio (SNR). For example, in the case of fading multipath channels, a MIMO configuration may increase system capacity by nearly M additional bits/cycle for each 3-dB increase in SNR.\nHowever, the widespread deployment of multi-antenna systems in wireless communications, particularly in wireless handset devices, has been limited by the increased cost that results from increased size, complexity, and power consumption. Providing separate RF chain for each transmit and receive antenna is a direct factor that increases the cost of multi-antenna systems. Each RF chain generally comprises a low noise amplifier (LNA), a filter, a downconverter, and an analog-to-digital converter (A/D). In certain existing single-antenna wireless receivers, the single required RF chain may account for over 30% of the receiver's total cost. It is therefore apparent that as the number of transmit and receive antennas increases, the system complexity, power consumption, and overall cost may increase. This poses problems for mobile system designs and applications.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "The outstanding sulfidation-resistant alloy available in the art has been alloy 6B invented by E. Haynes (U.S. Pat. No. 1,057,423) and marketed under the registered trademark STELLITE. STELLITE.RTM. alloy 6B is cobalt base and contains about 30% chromium, 4% tungsten, 1.1% carbon and is essentially free of iron and nickel.\nThe high cost and strategic limitations of cobalt prevent the full marketing of the alloy for wide spread use in combating sulfidation damage. The production costs of alloy 6B are especially high because of the difficulty in forging and hot and cold rolling this alloy. Furthermore, it is difficult to fabricate the alloy into components such as heat exchangers for applications.\nU.S. Pat. Nos. 4,195,987 and 4,272,289 disclose alloys containing iron, nickel, cobalt, chromium and selected metals including lanthanum to increase resistance to high temperature oxidation. A commercial alloy, marketed under the registered trademark HAYNES.RTM. alloy 556, is a typical example of this prior art. The alloy normally contains essentially about 18% cobalt, 22% chromium, 3% molybdenum, 2.5% tungsten, 20% nickel, 0.6% tantalum, 0.02% lanthanum and the balance iron with minor contents of nitrogen, manganese, aluminum, carbon and zirconium.\nU.S. Pat. No. 3,418,111 discloses Haynes alloy 188, well-known in the art for its resistance to high temperature oxidation. The alloy normally contains about 22% nickel, about 22% chromium, about 14% tungsten, 0.10% carbon, 0.03% lanthanum, and the balance essentially cobalt (about 40%).\nKnown in the art is UMCo-50 alloy or HAYNES alloy 150. The alloy contains normally about 28% chromium, about 50% cobalt and the balance iron with minor contents of carbon, manganese, and silicon. The alloy has good high temperature properties including stress-rupture and sulfidation resistance.\nMany prior art alloys, including those mentioned above, are used as components in industrial installations where resistance to chemical reactions such as oxidation and sulfidation is required. Equally the weldability and thermal stability characteristics must be acceptable.\nEach of the prior art alloys provides one or more of the desired characteristics but may be deficient in one or more of the other required characteristics. In some cases an alloy may provide nearly all the desired characteristics but its use may be limited because of the cost of raw materials and processing. Thus, the art is in need of an alloy that provides all of the desired characteristics at a lower cost."}
{"text": "Conventionally, in order to form structures on a substrate by using a process that minimizes manual machining, such process may involve techniques such as photolithography, electroplating, etching, and so forth. Photolithography has been a reliable method used to form structural patterns on a substrate, and a process of using photolithography during fabrication to pattern parts of a substance deposited on a substrate would typically involve one or more photo masks having a specific geometric pattern which would be transferred from a photo mask to a light-sensitive photoresist. After a period of exposure to light, chemical or high intensity photon treatments would be applied to remove a part of the light-sensitive photoresist to develop a patterned exposure into which a material could then be deposited on the material underneath the photo resist.\nHowever, such process would involve numerous expensive machines and masks as well as extremely clean operating conditions so that the overall costs of manufacturing would be increased. Also, in order to secure a photoresist onto a substrate, a polymeric resin such as epoxy would need to be sprayed on parts on a bonding surface before curing. Such process would involve additional chemicals and extra steps. Furthermore, even though electroplating has been a popular process used to deposit material into exposed parts of a photo resist on the material underneath the photo resist, conventional electroplating would also incur costs because of the use of electroplating equipment, electrodes, plating solutions, and etc. . . . .\nTherefore, there could be a more efficient and cost effective way of developing patterned structures on a substrate than above mentioned conventional practices."}
{"text": "1. Technical Field of the Invention\nThe present invention relates generally to pressure sensors and specifically to pressure sensors to be used in high pressure, cryogenic environments.\n2. Discussion of the Related Art\nPressure sensors are used to monitor fluid and gas pressures in a wide variety of applications. Many of these applications involve placing the sensors in environments that may damage the sensors or limit their accuracy. For example, in wind tunnel applications such as the National Transonic Facility at NASA Langley Research Center temperatures may be as low as approximately -173.degree. C. In the Space Shuttle Main Engine pressure must be sensed in the fuel supply lines. In this application gaseous and liquid oxygen or hydrogen are present at very high pressures and very low temperatures. The sensors must be able to operate within a temperature range from -253.degree. C. to 60.degree. C. and pressures from 0 to 5,000 psi. Additionally, they will be subject to 80 g vibrations from 25 to 2,000 Hz and up to 400 g impulse shock. Chemical resistance to O.sub.2 and H.sub.2 is also important to long term sensor survivability and reliability.\nBrown, et. al. (U.S. Pat. No. 5,454,270) disclose a hermetically sealed pressure sensor for use in a hostile environment. Generally, the device of Brown is for use in measuring fluid or gas pressures where the fluids or gases may damage the sensing device. Examples given include petrochemicals, freons, solvents, and alcohols. To protect the device from the corrosive effects of the hostile environment it is sealed in a plastic housing and only one face of the pressure transducer is exposed to the hostile environment. The '270 patent further discloses a stress isolation base to which a differential pressure transducer is attached. It is specified that the base is made of a ceramic material such as alumina, or other material having a similar coefficient of thermal expansion as silicon, the material from which the differential pressure transducer is made. This choice helps the pressure sensor to be accurate over a range of temperatures given in the disclosure to be approximately -40.degree. C. to +150.degree. C.\nMaurer (U.S. Pat. No. 5,351,550) discloses a pressure sensor comprising a pressure transducer, a housing member and a pressure sensor die having a diaphragm with at least one piezoresistive component disposed thereon.\nMaurer (U.S. Pat. No. 5,327,785) discloses a pressure sensor with an elastomeric member for heat dissipation.\nKurtz et. al. (U.S. Pat. No. 5,303,594) disclose a pressure transducer using polycrystaline diamond film. The advantages disclosed include high temperature sensing beyond the range available with silicon pressure sensors and improved output signal strength over silicon carbide sensors.\nChapman (U.S. Pat. No. 5,116,331) discloses a pressure transducer for use in cryogenic environments. The '331 patent discloses that by increasing boron dopant density in the piezoresistive bridge elements of a sensor from approximately 10.sup.16 boron/cm.sup.3 to &gt;1.3.times.10.sup.19 boron/cm.sup.3 the sensor becomes more thermally stable. Also disclosed are the drawbacks to highly doped sensors including propensity to mechanical failure and reduced pressure sensitivity.\nIn the '331 patent, a plurality of highly doped (10.sup.19 -10.sup.21 boron/cm.sup.3) silicon piezoresistive pressure sensors are mounted on a substrate for sensing pressures in a wind tunnel environment. Each pressure sensor is paired with a temperature sensor to provide for temperature correction to the sensors output in real time. Increased amplification is used to make up for the problem of reduced pressure sensitivity of the highly doped sensors. The sensor is mounted to a borosilicate glass substrate such as Corning, Inc.'s Pyrex 7740. Borosilicate is chosen to provide a coefficient of thermal expansion similar to that of highly doped silicon. In the '331 patent this is given as 2.5 ppm/C for highly doped silicon, and 3.2 ppm/C for Pyrex 7740, compared to 6.5 ppm/C for alumina, the material disclosed in the '270 patent to Brown above.\nSahagen (U.S. Pat. No. 5,088,329) discloses a sapphire force collector diaphragm having piezoresistive silicon films formed thereon. The piezoresistive films are arranged to form a Wheatstone bridge. One side of the force collector is in contact with the media being measured, the other, having the piezoresistive silicon films is not, thereby allowing the device to be used in high temperature or corrosive applications. The '329 patent notes that in a diaphragm type sensor there is a preferred region of the diaphragm in which the piezoresistive elements should be placed. Within a region having radius R1, corresponding in the '329 patent to that region of the diaphragm which is unsupported, there is a second region having radius R2, within which deflection of the diaphragm does not cause measurable stresses. Thus, the piezoresistive elements are preferably placed in the annular region between R1 and R2. The '329 patent discloses that R2 is preferably approximately 0.66 R1.\nGraeger, et. al. (U.S. Pat. No. 5,024,097) disclose a silicon body having piezoresistive elements formed thereon. The piezoresistive elements are further arranged to form a Wheatstone bridge. The silicon body has a blind hole forming a cavity between the silicon body and its substrate forming a diaphragm."}
{"text": "The present invention relates generally to granulating apparatus and, more specifically, to such apparatus of the wet aggitative type adapted for forming powder materials, particularly advanced ceramic and similar valuable materials, into granules utilizing an agglomerating liquid.\nIn various and diverse technologies, raw material processing conventionally requires the initial processing of raw material in powder form into larger agglomerated granules which may be more easily handled. Preferably, such agglomerated granules should be substantially spherical in shape and as uniform as possible in composition and density for ease of flowability and handling of the granules and to enhance the quality of end products formed therefrom. The formation of such uniform granules is particularly important in the preliminary processing of materials utilized in the emerging technological field of advanced ceramics. By way of example, such materials include alumina, silicon carbide, silicon nitride, barium titanate, various metal oxides, cermets such as tungsten carbide, partially stabilized zirconia, and cubic boron nitride. Advanced ceramic materials such as the foregoing find such highly technological applications as in the ferrite cores utilized in the deflection yokes for cathode ray tubes, adiabatic diesel engines and ceramic turbine engines, ceramic gas sensing devices, certain catalyst carriers and substrates, and various engineering ceramic applications such as high speed cutting tools, grinding media and the like. Such materials are conventionally formed into end products by a dry pressing or compacting operation. However, as initially processed, such materials are in a powdered form which characteristically is light, fluffy and therefore unmanageable and unsuitable for pressing operations. Accordingly, it is conventionally necessary to first form the powdered material into uniform granules which may be more readily handled and processed. Such advanced ceramic materials also are relatively expensive and valuable and, for the specialized and highly technological applications to which they are typically put, must be of optimum purity to insure a high degree of quality in the end product. Accordingly, the aforementioned requirements of uniform granule shape, composition and density are particularly acute and important in this field.\nVarious forms of granulating apparatus which have been developed over past years have proved entirely unsuitable and unacceptable for formation of granules of advanced ceramic materials. So-called spray dryers are currently employed widely for the mass production of advanced ceramic granules. In spray dryers, a slurry composed of powdered raw material and a binder or other agglomerating liquid is continuously sprayed in essentially discrete droplet or globular form continuously in a downward direction into the top end of an upright hollow vessel while heated drying air is directed upwardly within the vessel to drive off the agglomerating liquid thereby transforming the globules into essentially dry granules of the powdered material. While this apparatus produces generally uniform spherical granules having acceptable flowability characteristics and providing acceptable suitability for dry pressing operations, it is difficult to produce relatively dense granules due to the heat-induced evaporation of the agglomerating liquid which tends to cause outward migration of the finer particles of the powdered material toward the granule surface and thereby often produces hollow granules. Additionally, the operation of the apparatus is highly energy intensive, requiring approximately a 50% moisture level in the slurry for acceptable granulation and accordingly requiring a considerable amount of drying energy. As a result, the ratio of granule output to energy consumed is unacceptably low. The substantial drying requirements of such apparatus generally also require that the apparatus be undesirably large and bulky. Furthermore, since the granulation process performed by such apparatus is a continuous one, the apparatus is basically unsuitable for use in forming small quantities of granules. When it is desired to change the powdered material being granulated, it is of course necessary to clean the entire apparatus which is difficult and time consuming to do and often in any event results in cross contamination from one material to the next.\nExtruding machinery is also available for forming moistened powder material into generally cylindrical pellet-type form by the forced extrusion of the material through a die plate or screen. Disadvantageously, it is difficult to produce pelletized granules by this operation to smaller than approximately a 0.7 mm. diameter due to rapid resultant wearing of the screen or die plate and concomitant contamination of the pelletized granules. Moreover, the cylindrical shape of the pellets significantly restricts their flowability and, as a result, such pellets are often further processed in a spheronizing vessel wherein the pellets are repetitively beaten against the vessel walls by a rotating plate to deform and plasticize the pellets into spherical form.\nDisintegration-type granulating machinery is also available for breaking large agglomerations of material into smaller pieces, but such apparatus produce generally unacceptable irregular and poorly flowable granule shapes and also suffer significant losses of relatively fine particles resulting in a low granule yield rate.\nFluidized bed granulating machines have been developed which essentially combine the functions of an extruder, a spheronizer and a dryer for granule formation. In such apparatus, powdered raw material is fed into a cylindrical vessel having a stationary screen or plate with openings therethrough, the powdered material being treated with an agglomerating liquid spray on the charge side of the screen or plate while heated drying air is forced through the screen or plate openings from the opposite side thereof to form the powder into dry granules. Disadvantageously, such granulating apparatus produces a widely varying range of granule shapes and sizes which inhibit good flowability of the granules. The apparatus is often subject to rapid component wearing with resultant granule contamination, and also is not susceptible of accurate control and reproducibility of granule size.\nIn another form of aggitative-type granulating apparatus, a charge of powdered material is processed in batch form in a vessel having a rotating outer annular mixing or stirring rotor and an inner impulse rotor therewithin for compatibility circulating and mixing the powdered charge material, with an agglomerating liquid spray being provided for forming the powder material into granules as circulation progresses. To enhance the powdered material distribution and granule formation, blade-like implements are typically provided on each rotor and the rotors may be rotated in opposite directions with the central impulse rotor being operated at a greater speed than the outer stirring rotor. Examples of this type of granulating apparatus are disclosed in Japanese Patent Publication Nos. 58-12050, 59-21649, 59-55338, and 59-59239, commonly owned with the present invention. While these apparatus are acceptably operative for granule formation, such apparatus produces widely ranging granule sizes with poor capability for accurate repeatable control of granule size and such apparatus also suffers rapid wearing of its internal operating components with attendant contamination of the granules formed, all of which makes this apparatus generally unacceptable for forming granules of advanced ceramic and other valuable materials.\nIt is accordingly an object of the present invention to provide an improved granulating apparatus of the aggitative type last above-described, which is capable of repeatably accurate production of spherical granules of uniform density and composition and high purity without significant contamination, thereby being uniquely suitable for advanced ceramic granule formation."}
{"text": "1. Field of the Invention\nThe invention concerns a switching element for a valve train of an internal combustion engine, preferably for valve deactivation, having an outer part and an inner element that is axially displaceable in a bore therein, wherein the outer part and the inner element each have at least one receptacle aligned with each other in an axially separated relative position achieved by a lost motion spring, wherein one or more movable coupling means for coupling the inner element with the outer part are applied in one of the receptacles toward the other receptacle in their positions relative to one another, wherein a first upper stop for defining the relative position is applied between the inner element and the outer part, and wherein a hydraulic lash adjuster with a pressure piston is installed in the inner element, which pressure piston is fixed against moving axially out of the inner element by a second upper stop.\n2. Description of the Related Art\nA switching element of this type is disclosed in DE 199 15 531 that is considered generic. This switching element is shown as a switchable cam follower for a tappet push rod drive. An upper stop for defining the relative position is realized through a piston-like element arranged in the inner element. This piston-like element projects radially outward into a longitudinal groove of the outer part. In the axially extended state of the inner element relative to the outer part, the piston-like element contacts an end of the longitudinal groove. The aim of this is to achieve an aligned positioning of a coupling bore provided in the outer part and a piston arranged in the inner element for coupling.\nA drawback of this prior art is that adjustment of lash in the coupling is relatively complicated and expensive. It is clear that the receptacle in the outer part (here a coupling bore) for receiving the piston must be designed with a slight lash relative to the outer peripheral surface of the piston. This lash and a vertical position vary from one switching element to the other depending on the manufacturing conditions. The relatively broad range of variation of this mechanical free travel in the switching elements is, however, not desirable.\nTherefore, to adjust the coupling lash or keep its variance within an acceptable range, the pistons are classified for locking purposes in groups. This is extremely complicated and expensive from the manufacturing and measuring point of view. For example, switching elements must be completely assembled, the lash then measured, following which the switching element must again be disassembled and mated to a suitable coupling piston. It is equally conceivable to classify the upper stops on the longitudinal groove of the outer element.\nAnother upper stop is provided in the aforesaid prior art for a pressure piston of the lash adjuster and is configured as a ring.\nIf two pistons are provided for coupling, as is the case in DE 4,206,166, the aforesaid stop measures prove to be nearly impossible. The aligned position of the coupling bores situated diametrically opposite each other in the inner element is realized when the two axially movable parts of the switching element make contact with the base circle of the cam. An adjustment of the coupling lash in this case is accomplished by extremely complicated manufacturing and measuring techniques by pairing suitable switching elements (in this case, cup tappets) with cam pairs or camshafts. Under certain circumstances, it is necessary to tolerate an excessive lash variation.\nThus, the object of the invention is to provide a switching element of the aforementioned type in which the stated drawbacks are eliminated by simple means."}
{"text": "With the advent of pagers and mobile phones the wireless service industry has grown into a multi-billion dollar industry. The Cellular Telecommunications and Internet Association calculates that 120 million Americans own a mobile telephone—about half of the U.S. population. As the development and availability of mobile telephones progresses the benefits of mobile telephones are reaching more and more people. The inclusion of games on mobile telephones has also enjoyed increasing popularity. It is common for modem cellular phones and Personal Digital Assistant—based communication devices to include features that are used for social purposes, such as messaging and cameras. While there are many types of games available for these phones, they do not take full advantage of these social capabilities.\nIn addition, revenues from massively multiplayer games, online games played between large numbers of participants, have increased as they have gained popularity. Massively multiplayer games, however, require users to be located at desktop computers or other static locations because the games require use of a computer. This leads many massively multiplayer game users to become uncomfortable or restless after long periods of time. Users of these games desire more freedom to do other things and have access to other locations as they play massively multiplayer games.\nTherefore a need exists to overcome the problems with the prior art as discussed above."}
{"text": "1. Technical Field\nThe present disclosure relates to systems, devices and methods for performing a medical procedure. More particularly, the present disclosure relates to heat-distribution indicators and thermal zone indicators suitable for use during thermal ablation, electrosurgical systems including the same, and methods of directing energy to tissue using the same.\n2. Discussion of Related Art\nElectrosurgery is the application of electricity and/or electromagnetic energy to cut, dissect, ablate, coagulate, cauterize, seal or otherwise treat biological tissue during a surgical procedure. When electrical energy and/or electromagnetic energy is introduced to tissue, the energy-tissue interaction produces excitation of molecules, creating molecular motion that results in the generation of heat. Electrosurgery is typically performed using a handpiece including a surgical instrument (e.g., end effector, ablation probe, or electrode) adapted to transmit energy to a tissue site during electrosurgical procedures, an electrosurgical generator operable to output energy, and a cable assembly operatively connecting the surgical instrument to the generator.\nTreatment of certain diseases requires the destruction of malignant tissue growths, e.g., tumors. Electromagnetic radiation can be used to heat and destroy tumor cells. Treatment may involve inserting ablation probes into tissues where cancerous tumors have been identified. Once the probes are positioned, electromagnetic energy is passed through the probes into surrounding tissue. The application of ultrasound imaging is one of the cost-effective methods often used for tumor localization and ablation device placement.\nThere are a number of different types of apparatus that can be used to perform ablation procedures. Typically, apparatus for use in ablation procedures include a power generating source, e.g., a microwave or radio frequency (RF) electrosurgical generator, that functions as an energy source, and a surgical instrument (e.g., microwave ablation probe having an antenna assembly) for directing the energy to the target tissue. The generator and surgical instrument are typically operatively coupled by a cable assembly having a plurality of conductors for transmitting energy from the generator to the instrument, and for communicating control, feedback and identification signals between the instrument and the generator.\nUsing electrosurgical instruments to ablate, seal, cauterize, coagulate, and/or desiccate tissue may result in some degree of thermal injury to surrounding tissue. For example, electrosurgical desiccation may result in undesirable tissue damage due to thermal effects, wherein otherwise healthy tissue surrounding the tissue to which the electrosurgical energy is being applied is thermally damaged by an effect known in the art as “thermal spread”. During the occurrence of thermal spread, excess heat from the operative site can be directly conducted to the adjacent tissue and/or the release of steam from the tissue being treated at the operative site can result in damage to the surrounding tissue. The duration of the activation of the generator is directly related to the heat produced in the tissue. The greater the heat produced, the more the potential for thermal spread to adjacent tissues.\nCurrently available systems and methods for controlling an electrosurgical generator during electrosurgery may include a clinician monitoring and adjusting, as necessary, the amount of energy delivered to a tissue site through current, voltage, impedance, and/or power measurements such that an appropriate tissue effect can be achieved at the tissue site with minimal collateral damage resulting to adjacent tissue. These systems and/or methods typically require a clinician to translate the desired tissue effect to a power setting on an electrosurgical generator and, if necessary, adjust the power setting to compensate for tissue transformations (e.g., desiccation of tissue) associated with the electrosurgical procedure such that a desired tissue effect may be achieved.\nAs can be appreciated, limiting the possibility of thermal spread or the like during an electrosurgical procedure reduces the likelihood of unintentional and/or undesirable collateral damage to surrounding tissue structures which may be adjacent to an intended treatment site. Controlling and/or monitoring the depth of thermal spread during an electrosurgical procedure may aid a clinician in assessing tissue modification and/or transformation during the electrosurgical procedure.\nMedical imaging has become a significant component in the clinical setting and in basic physiology and biology research, e.g., due to enhanced spatial resolution, accuracy and contrast mechanisms that have been made widely available. Medical imaging now incorporates a wide variety of modalities that noninvasively capture the structure and function of the human body. Such images are acquired and used in many different ways including medical images for diagnosis, staging and therapeutic management of malignant disease.\nBecause of their anatomic detail, computed tomography (CT) and magnetic resonance imaging (MRI) are suitable for, among other things, evaluating the proximity of tumors to local structures. CT and MRI scans produce two-dimensional (2-D) axial images, or slices, of the body that may be viewed sequentially by radiologists who visualize or extrapolate from these views actual three-dimensional (3-D) anatomy.\nMeasurements and quantitative analysis for parameters such as area, perimeter, volume and length may be obtained when object boundaries are defined. A boundary in an image is a contour that represents the change from one object or surface to another. Image segmentation involves finding salient regions and their boundaries. A number of image segmentation methods have been developed using fully automatic or semi-automatic approaches for medical imaging and other applications. Medical image segmentation refers to the delineation of anatomical structures and other regions of interest in medical images for assisting clinicians in evaluating medical imagery or in recognizing abnormal findings in a medical image. Structures of interest may include organs or parts thereof, such as cardiac ventricles or kidneys, abnormalities such as tumors and cysts, as well as other structures such as bones and vessels.\nMedical image processing, analysis and visualization play an increasingly significant role in disease diagnosis and monitoring as well as, among other things, surgical planning and monitoring of therapeutic procedures. Unfortunately, tissue heating and thermal damage does not create adequate contrast in ultrasound images to allow determination of the size of an ablated zone and assessment of the margins of ablated tissue."}
{"text": "1. Field of the Invention\nThe present invention relates to the preparation of semiconductor devices, and more particularly to a method used to establish an electrical field near a surface of a semiconductor device, and in particular, a charge-coupled imaging device.\n2. Description of the Prior Art\nCharge-coupled devices are typically made of silicon and are used as solid-state imagers by taking advantage of the properties of a silicon crystal lattice. In the crystalline form, each atom of silicon is covalently bonded to its neighbor. Energy greater than the energy gap of about 1.1 V is required to break a bond and create an electron hole pair. Incident electromagnetic radiation in the form of photons of wavelength shorter than 1 um can break the bonds and generate electron hole pairs.\nThe wavelength of incoming light and the photon absorption depth are directly related, the shorter the wavelength, the shorter the penetration depth into the silicon. Silicon becomes transparent at a wavelength of approximately 1100 nm and is essentially opaque to light at wavelengths shorter than 400 nm. High energy particles, X-rays and cosmic rays can break many thousands of bonds; therefore, excessive exposure to these can cause damage to the crystal lattice. Bonds can also be broken by thermal agitation. At room temperature, approximately 50 bonds per second per um.sup.3 are broken and recombined on a continuous basis. The rate of electron hole pair generation due to thermal energy is highly temperature-dependent and can be reduced arbitrarily through cooling.\nIn order to measure the electronic charge produced by incident photons, it was required to provide a means for collecting this charge. Thus, the potential well concept was developed, wherein a thin layer of silicon dioxide is grown on a section of silicon, and a conductive gate structure is formed over the oxide. The gate structure is formed in an array of columns and rows, thus making it possible, by applying a positive electrical potential to various gate elements, to create depletion regions where free electrons generated by the incoming photons can be stored.\nBy controlling the electrical potential applied to adjacent gates, the depletion region, or well, containing the free electrons can be caused to migrate along a column or row, so that the signal may eventually be output at the edge of the array.\nTypically, the gate structure is arranged with multiple phases, particularly three phases, so that the potential wells may be easily migrated through the silicon to an output device.\nIn reality, the wells and the migration of the wells is not carried out along the surface of the silicon-silicon dioxide interface, but takes place in a buried channel below the surface. The buried channel is free of interference from interface states and thus assures effective charge transfer from well to well. The operation of a charge-coupled device is somewhat analogous to that of a bucket brigade circuit commonly used to delay electrical signals.\nBecause the charge from the wells located far from an output amplifier must undergo hundreds of transfers, the charge transfer efficiency of a charge-coupled device is most important, as is the quantum efficiency and the spectral response. These considerations are particularly important when extremely low light levels are to be sensed.\nLight normally enters the charge-coupled device by passing through the gates in the silicon dioxide layer. The gates are usually made of very thin polysilicon, which is reasonably transparent to long wavelengths but becomes opaque at wavelengths shorter than 400 nm. Thus, at short wavelengths, the gate structure attenuates incoming light.\nIn an effort to overcome this difficulty, it has become the practice to uniformly thin a charge-coupled device to a thickness of approximately 10 um, using acid etching techniques. Using a thinned charge-coupled device, it then becomes possible to focus an image on the backside of the charge-coupled device, where there is no gate structure that will attenuate the incoming light. Thinned charge-coupled devices exhibit high sensitivity to light from the soft X-ray to the near-infrared region of the spectrum.\nFIG. 1A illustrates schematically a cross-section of a typical thick-bodied charge-coupled device. The device includes a silicon body 2, a silicon dioxide layer 4 and a gate array 6 formed on the silicon dioxide layer. Incoming light is illustrated by arrows 8 as illuminating a front side of the silicon 2. FIG. 1B illustrates a cross-section of a thinned charge-coupled device with light illuminating a backside. The thinned charge-coupled device, having a thickness of approximately 10 um, has improved quantum efficiency and UV spectral response.\nIn order to improve the quantum efficiency of the charge-coupled device, it is desirable that the free electrons accumulate in the buried channel near the front surface of the charge-coupled device under the gate structure. It has become the practice to develop an electronic field near the rear surface of the charge-coupled device to drive electrons towards the front pixel surface, so that the free electrons may be concentrated in the buried channel. This practice is also used with thinned charge-coupled devices and in such cases, the rear surface of the device is subjected to ion implantation of, for example, boron, followed by laser annealing for full implant activation.\nOne difficulty that has been encountered with this practice is that the boron concentration maximum is not at the rear surface of the charge-coupled device, but rather is buried below the surface at a depth ranging from 2000.ANG. to over 5000.ANG.. As a result, the rear surface of the silicon is at a different potential than the sub-surface silicon. This situation may not be undesirable for front-illuminated devices or for rear-illuminated devices for light in the visible spectrum, where photon penetration is greater than 5000.ANG.. However, for rear-illuminated devices, as wavelengths become shorter, the penetration depths shorten, and the photo-electrons see a potential barrier and are essentially driven to the rear surface, as opposed to the front surface. As a result, the quantum efficiency of the device becomes exceedingly low.\nTo reduce this undesirable phenomenon, ion implantation is generally carried out at low ion beam energies equal to or slightly greater than 5 keV. Such low ion beam energies result in an implantation concentration maximum nearer to the rear surface as, for example, at approximately 1500.ANG.. However, with low implantation energies, the ion distribution profile tends to be rather broad, and a good sharp implantation distribution is not realized. Further broadening of the distribution profile occurs when the implantation is laser activated into the silicon matrix.\nUsing such techniques, the sub-surface maximum is satisfactory for visible light, but becomes a major concern and results in loss of quantum efficiency as wavelengths become shorter and the photon penetration depths fall below 2000.ANG.. When this occurs, the photo-electron sees a potential wall that results in reduced quantum efficiency.\nThe problem of generating an electrical field at a surface of semiconductor material is encountered in many other fields of endeavor. The problem is not limited to the preparation of charge-coupled imaging devices. In the preparation of solar cells, the same difficulty is realized.\nIn order to more fully understand the techniques in ion implantation and the ion distribution profiles that may be generated, one may refer to the text entitled VLSI Technology, edited by S. M. Sze, and published by McGraw-Hill Book Company, and in particular, chapter 6 of said text."}
{"text": "1. Field of the Invention\nThis invention relates to an automotive window retention system for motor vehicle windows which are fixed (i.e., stationary) with respect to the vehicle body. Specifically, the invention relates to an automotive window retention system for mounting a window assembly onto a recessed generally L-shaped peripheral flange of a window aperture formed in an automotive vehicle body which employs both hook and loop fabric fasteners and a curable adhesive embedded therein to provide a very secure yet facile assembly of the window to the vehicle body.\n2. Description of the Prior Art\nVarious automotive window retention systems are well known, one of which is represented in FIG. 1 of the drawings. As shown therein, a fixed automotive window opening normally consists of a metal interior body panel generally designated 10 welded to a metal exterior body panel generally designated 12, the latter of which has a peripherally recessed L-shaped flange, generally designated 14, peripherally extending about the desired automotive window opening. Flange 14 consists of a lower leg 15 perpendicular to the exterior body panel 12 and an upper leg 16 abutting an upper edge section 19 of the interior body panel 10 and disposed generally normally to the lower leg 15. Mounted on the lower leg 15 is one or more rubber or plastic spacers 17 and mounted on the upper leg 16 is one or more rubber or plastic semispherical spacers 18 which serve to properly position and space the glass window 20 in the recessed body panel opening defined by the L-shaped flange. A polyurethane adhesive is applied into the space between the upper leg 16 and the interior side of the window glass panel 20 which, upon curing, serves to permanently secure the window in the body panel opening.\nAn interior window trim 22 is secured to the interior side of the interior body panel via Velcro.RTM.-type hook and loop fabric fastening elements 23, 24, one of which is secured to the interior trim 22 and the other of which is secured to the interior side of the upper edge section 19 of the interior body panel 10. An outer window trim 26 covers the outer edge of the window 20 and a portion of the exterior body panel 12 and is secured to the latter via a screw 27 extending through the exterior window trim 26 and the body panel 12 properly positioned by a spacer 28.\nAs can be appreciated from the above description, this prior art retention system requires numerous parts and assembly steps, as a result of which it is quite costly in terms of both actual material costs and in labor. Furthermore, so far as is known, no presently available window retention system is available which minimizes the mounting materials used, saves steps and time, and allows for an easy and facile mounting and retention of the window system in a manner comparable to that according to the present invention."}
{"text": "Embodiments generally relate to web based security management. More particularly, embodiments relate to the intelligent aggregation of platform device interaction information associated with web applications.\nEmerging markup languages such as HTML5 (Hypertext Markup Language 5, e.g., HTML5 Editor's Draft 8 May 2012, W3C), LLVM (e.g., LLVM 3.1, May 22, 2012, llvm*org), and other runtime or just in time (JIT) environment languages may support more robust multimedia related web platform development. The use of these advanced languages by an application developer, however, may also expose end user platform components such as graphics processors, memory, sensors, and so forth, to web applications, wherein the exposure of such components may lead to security concerns."}
{"text": "The invention relates to a process for brazing metallic materials, an aqueous flux suspension and a novel flux for use in a brazing process and a preparation process for the novel flux.\nDuring brazing, metal materials are joined at a temperature of above 500.degree. C. with the aid of a molten additional metal (solder). The melting point of the solder is lower than the melting point of the materials, so that the solder joints of many metals can be separated in a nondestructive manner by melting the solder again.\nDuring soldering, oxides and other interfering top layers on the metal surface present a problem: the surface must be metallically pure so that a satisfactory solder joint forms. This is achieved using fluxes, which are generally brushed on, sprayed on or applied as a solder covering on the material.\nFluxes based on alkali metal salts, preferably potassium salts, of complex aluminum fluorides are particularly suitable, for example for brazing light metal materials, in particular alloyed (for example with magnesium) or unalloyed aluminum materials, since they are non-corrosive and non-hygroscopic. Fluxes of this type which are, for example, mixtures of KAlF.sub.4 and K.sub.3 AlF.sub.6 or mixtures of KAlF.sub.4 and K.sub.2 AlF.sub.5 are already known, and the last mentioned compound can optionally be present as the hydrate.\nThe surface-cleaning action of a flux is evident from, inter alia, the fact that, after melting, the solder runs on the surface of the component or of the components. The better the surface of the component has been cleaned by the flux, the more readily the solder runs on that surface."}
{"text": "This invention relates to carbon-carbon composite materials. In one aspect this invention relates to carbon-carbon composite materials which are resistant to oxidation. In another aspect, this invention relates to a method for providing oxidation resistance for carbon-carbon composite materials.\nReinforced composites are used in a wide variety of applications. The best known composites are made from two-dimensional fabrics and/or fibers dispersed in a resin or plastic matrix. These composites are basically a resin or plastic structure to which reinforcing fabrics or fibers have been added to enhance the physical properties of the structure.\nAdvances in the field of aerospace technology have created a need for high strength, temperature-resistant materials. For many applications, this need is satisfied by carbon-carbon composite materials. These materials utilize a carbon matrix, as opposed to a resin or plastic matrix.\nA wide range of multidirectional reinforced composite structures are now available. The simplest of these structures is obtained by stacking unidirectional fibers or sheets with alternating layers oriented in different directions, or by stacking woven sheets. More complex structures provide three-dimensional reinforcement. The simplest of the three-dimensional structures is the three-directional (3D) structures which generally has reinforcing elements which are mutually orthogonal. The most complex three-dimensional structure is a thirteen-directional (13D) structure. The thirteen directions, with reference to a cube, form three subgroups; the three edges, the four long diagonals, and the six diagonals of the faces.\nThe reinforced carbon-carbon composite structures are fabricated from graphite or carbon yarn or rods. The term yarn includes continuous filament yarns and yarns spun from short fibers, and comprises a plurality of filaments or fibers combined to make up a desired end count. Rods are produced by a pultrusion process whereby unidirectional groups of carbon or graphite yarn are assembled and impregnated with a thermosetting or thermoplastic resin or binder. The impregnated yarn groups are drawn through a die which is warmed to a desired temperature and which has a suitable shape.\nThe carbon or graphite yarns or rods are assembled into the desired geometric structure. If desired, the yarn may be impregnated with a suitable resin or binder prior to assembly.\nThe composite is formed either by sintering the reinforcement structure by solidifying the impregnated precursor, thereby avoiding the requirement for other materials, or by the dry or the liquid process, or by a combination of these methods. The dry process consists of depositing pyrolytic carbon inside the structure of the reinforcement by decomposition of a hydrocarbon gas such as methane. In the liquid process, the porous texture of the reinforcement is impregnated with a thermosetting resin or a thermoplastic carbon precursor, such as a phenolic resin, a furanyl resin, coal tar pitch, or the like, that is converted to carbon by heat treatment. Following carbonization, the structure is graphitized. The impregnation, carbonization, graphitization cycle is repeated as often as necessary to densify the composite to a desired degree.\nThe process of densification of the composite generally comprises heat treatment at a temperature in the range of 2500.degree. to 3000.degree. C. and may include isostatic pressing at pressures up to about 15,000 psi in an oxygen-free environment.\nMany applications for carbon-carbon composites have been proposed or implemented. The use of such composites for re-entry heat shield applications has been demonstrated. Ehrenreich, U.S. Pat. No. 3,672,936, discloses disk brake pads made of such composites. The use of these materials for turbine disk and blade components, for propulsion system nozzles, thrust chambers, and ramjet combustion liners, and for re-entry vehicle nosetip applications has been investigated.\nIn the presence of an inert atmosphere, carbon has a sublimation point in excess of 3500.degree. C. When heated in excess oxygen, carbon burns at about 400.degree. C. For certain applications an allowable amount of wasting away due to combustion can be designed into the structure. For example, a re-entry vehicle heat shield is intended for a useful life of one cycle. For other applications, use and multiple re-use may be desired, in which case wasting away is to be avoided.\nOxidation resistance can be provided for carbon or graphite materials by depositing silicon over the carbon. Rubisch, U.S. Pat. No. 3,553,010, discloses that flame injection applied silicon reacts with the carbon of the underlying body forming silicon carbide when the operational temperature of the protected parts exceeds about 550.degree. C., which leads to a protective layer of silicon carbide which exhibits oxidation resistance at a relatively high temperature. When heated in the presence of oxygen, the silicon carbide is converted to silicon dioxide which has a softening temperature of about 1700.degree. C.\nThe protective layer of silicon carbide is highly flaw sensitive. If this protective layer is breached, the underlying carbon fibers can be quickly oxidized at temperatures below 1700.degree. C. due to the inability of the silica layer to flow and seal the ends of the damaged fibers. Breach of the silicon carbide layer may occur during use of the composite structure, as for example during high stress maneuvering periods of missiles, or it may occur during densification following application of the silicon carbide layer, as when pressure and/or temperature changes are inadvertently done too rapidly. Rapid thermal transients induced during component use can also be a major cause of cracking due to the thermal expansion difference between carbon and silicon carbide. Accordingly, what is desired is a protective coating capable of providing oxidation resistance for a carbon-carbon composite, which coating is also capable of flowing into and sealing a damaged area.\nAccordingly, it is an object of the present invention to provide an improved oxidation resistant coating for a carbon-carbon composite material.\nAnother object of the present invention is to provide a method for protecting the fiber bundles in a carbon-carbon composite against oxidation\nOther objects and advantages of the present invention will be apparent to those skilled in the art from a reading of the following detailed disclosure of the invention as well as the appended claims."}
{"text": "This invention relates generally to containers, particularly those that may be utilized for packaging, shipping and dispensing spools of wire, cable and the like.\nMore specifically, this invention relates to a combined container and dispenser reel that inexpensively provides means for both housing a spool of cable and the like during shipment thereof and housing said spool of cable and the like during actual dispensing and use thereof. This invention also relates to an inexpensively manufactured, combined packaging and dispensing means for a spool of cable and the like, in which said spool of cable may be mechanically engaged and locked into place within said container and dispensing device without a need for any special tools or components that would otherwise add to the cost of manufacturing."}
{"text": "1. Field of the Invention\nThe present invention relates to the use of a time-domain equalizer (TEQ) algorithm in a discrete multi-tone transceiver (DMT) and, more particularly, to a time-domain equalizer algorithm which operates as both a channel shortening filter and a noise whitening filter.\n2. Related Art\nThe fast, efficient and error-free transmission of digital information from one point to another has become increasingly important. Many communications systems exist which permit digital information to be transmitted over various types of communication channels, such as wireless channels, fiber-optic channels, and wire line channels.\nThe present invention will be described in the context of a wire line communications channel, such as a telephone line which utilizes a twisted pair of copper wires. It is noted that the use of the present invention is not limited to wire line systems as those skilled in the art will appreciate from the discussion hereinbelow.\nA modem is typically used to transmit and receive digital data over a telephone line. Modems employ a modulator to transmit the digital data over the telephone line and a demodulator to receive digital data from the telephone line. One common modulation technique is known as digital multi-tone modulation which requires a discrete multi-tone transmitter and a discrete multi-tone receiver at each modem in a communication system. Often, those skilled in the art refer to such modems as employing a DMT physical layer modulation technique.\nReference is now made to FIG. 1 which is a block diagram of a conventional DMT communications system 1. The system 1 includes a DMT transmitter 10, a transmission channel 20, and a DMT receiver 30. The DMT transmitter 10 includes a symbol generator 12, an inverse fast fourier transform (IFFT) modulator 14, and a cyclic prefix generator 16. The DMT transmitter 10 receives an input bit stream b(n) which is fed into the symbol generator 12. The symbol generator 12 produces a signal X(k) which is fed into the IFFT modulator 14. X(k) is a complex signal (i.e., a signal understood by those skilled in the art to comprise both a real and an imaginary component) formed by mapping pairs of bits of the input bit stream b(n) into a complex data space such that the complex signal X(k) has a length of N/2 samples. Symbol generator 12 also augments the signal X(k) with a complex conjugate to obtain a conjugate symmetric signal of N samples.\nThe IFFT modulator 14 performs an N-point inverse fast fourier transform on the conjugate complex signal X(k) to obtain the sampled real signal x(n). Since X(k) is a symmetric signal, the output of the IFFT modulator 16 is a real signal x(n). The real signal x(n) may be thought of as the summation of a plurality of cosine functions each having a finite length and a different frequency, phase, and amplitude, where these frequencies are multiples of a fundamental frequency. Since each of the cosine functions has a finite duration, x(n) is a varying amplitude discrete signal having a finite duration spanning N samples.\nFor the purpose of simplifying equations which will be discussed below, the transmission channel 20 is modeled as including a D/A converter 22, transmit filter (not shown), a receive filter (not shown), and an A/D converter 26 on either end of a wire loop 24. Those skilled in the art will appreciate that a practical system will employ the D/A converter 22 (and the transmit filter) in the DMT transmitter 10 and will employ the A/D converter 26 (and the receive filter) in the DMT receiver 30.\nThose skilled in the art will appreciate that the frequency spectrum of x(n) may be thought of as a plurality of orthogonal (SIN X)/(X) functions, each centered at a respective one of the frequencies of the cosine functions of x(n).\nx(n) is transmitted over the channel 20 to the DMT receiver 30. Since the transmission channel 20 has a non-ideal impulse response h(n), the received signal y(n) will not exactly match x(n). Instead, y(n) will be a function of the convolution of x(n) and h(n). Typically, h(n) will look substantially like the curve shown in FIG. 2. The non-ideal characteristic of h(n) introduces an amount of interference (specifically intersymbol and interchannel interference) which should be compensated for in both the DMT transmitter 10 and the DMT receiver 30.\nA common technique in compensating for the non-ideal impulse response of the transmission channel 20 is to introduce a so-called guard band at the beginning of each finite duration signal x(n) to produce xxe2x80x2(n). The cyclic prefix generator 16 performs this function. The guard band is typically formed of the last V samples of x(n) for each DMT symbol. If the length of the impulse response h(n) of the transmission channel 20 is less than or equal to V+1, then the guard band of length V will be sufficient to eliminate the interference cause by the impulse response h(n). The guard band is commonly referred to in the art as a xe2x80x9ccyclic prefixxe2x80x9d (CP).\nUnfortunately, the impulse response h(n) of a typical transmission channel 20 may be excessively long, requiring cyclic prefix lengths which substantially reduce the rate at which digital bits are transmitted across the transmission channel 20. The DMT receiver 30, therefore, employs signal processing techniques which effectively shorten the impulse response h(n) of the transmission channel 20, thereby permitting a corresponding reduction in the length of the cyclic prefix required at the DMT transmitter 10.\nThe DMT receiver 30 includes a time-domain equalizer (TEQ) 32, a window circuit 34, a fast fourier transform (FFT) demodulator 36, and a bit generator 38. The time-domain equalizer 32 is a finite impulse response (FIR) filter designed to compensate for the non-ideal impulse response h(n) of the transmission channel 20. In particular, the time-domain equalizer 32 employs a finite number of coefficients (T) which are calculated to compensate for the non-ideal impulse response of the transmission channel 20. As will be discussed in more detail below, the time domain equalizer 32 operates on the impulse response h(n) of the channel 20 such that the combined impulse response heff(n) of the channel 20 and the time domain equalizer 32 has maximum energy within a limited band of samples. This may be thought of as xe2x80x9cshorteningxe2x80x9d the effective impulse response of the channel 20. The output of the time domain equalizer is zxe2x80x2(n).\nA window circuit 34 is employed to remove the cyclic prefix from zxe2x80x2(n) to obtain z(n). The signal z(n) is input into the FFT demodulator 36 (which is understood to include a frequency domain equalizer function) to produce the complex symmetric signal X(k). After the complex conjugate portion of the signal X(k) is removed, the bit generator 38 maps the complex signal X(k) into an output bit stream b(n), which theoretically matches the input bit stream b(n).\nSeveral algorithms exist for calculating the T coefficients of the time-domain equalizer 32. One such algorithm is referred to as the least squares based pole zero cancellation algorithm (hereinafter xe2x80x9cPROCESS #1xe2x80x9d), which is discussed in detail in P. J. Melsa, R. Y. Younce and C. E. Rohrs, xe2x80x9cImpulse Response Shortening for Discrete Multitone Transceivers,xe2x80x9d IEEE Trans. On Comm. Vol. 44, No. 12, pp. 1662-71, December 1996, the entire disclosure of which is hereby incorporated by reference. Another such algorithm is referred to as the optimal shortening algorithm (hereinafter xe2x80x9cPROCESS #2xe2x80x9d), which is also discussed in detail in the above referenced IEEE publication. Still another algorithm is referred to as the eigenvector approach using the power method (hereinafter xe2x80x9cPROCESS #3xe2x80x9d), which is discussed in detail in M. Nafie and A. Gatherer, xe2x80x9cTime-Domain Equalizer Training for ADSL,xe2x80x9d Proc. ICC, pp. 1085-1089 (1997), the entire disclosure of which is hereby incorporated by reference.\nAlthough the above techniques for calculating the coefficients of the time-domain equalizer 32 address the issue of xe2x80x9cshorteningxe2x80x9d the effective impulse response of the transmission channel 20, they do not adequately address the problem of interference caused by noise. The system 1 is susceptible to injected white noise and colored noise. Colored noise exhibits concentrated spectral energy at some frequencies and relatively little spectral energy at other frequencies. White noise exhibits a substantially constant amount of spectral energy at all frequencies.\nWhile the conventional DMT receiver 30 of FIG. 1 operates optimally when white noise is present at the output of the time domain equalizer 32, it is susceptible to interchannel interference when colored noise is present. This is a particularly difficult problem when the colored noise exhibits spectral nulls.\nColored noise may be present at the output of the time domain equalizer 32 because (i) additive colored noise was injected into the signal xxe2x80x2(n) as it was transmitted over the transmission channel 20; and/or (ii) the time domain equalizer 32 itself introduces spectral shaping (especially spectral nulls) into the signal zxe2x80x2(n). Thus, even if the transmission channel 20 does not introduce additive colored noise into the received signal y(n), the time domain equalizer 32 may itself introduce spectral coloration into the additive noise of signal zxe2x80x2(n). Consequently, although the time domain equalizer 32 may produce a xe2x80x9cshorterxe2x80x9d effective impulse response heff(n), it may degrade system performance by introducing colored noise (especially spectral nulls) into zxe2x80x2(n). In particular, the rate at which data bits b(n) are transmitted over the transmission channel 20 and the error rate of such transmission may be adversely affected by colored noise at the output of the time domain equalizer 32.\nAccordingly, there is a need in the art for an improved DMT communication system which is capable of (i) compensating for the non-ideal impulse response of a transmission channel; (ii) compensating for additive colored noise introduced by the transmission channel; and/or (iii) mitigating against the spectral coloration of additive noise by the time domain equalizer.\nIn order to overcome the disadvantages of the prior art, an apparatus for receiving a discrete multi-tone signal over a communications channel having an impulse response h(n), energy of the impulse response being substantially concentrated in a first band of samples, the apparatus comprising:\na receiver for receiving the discrete multi-tone signal; and\na T coefficient finite impulse response time-domain equalizer included in the receiver, the time-domain equalizer having an output additive noise signal, the T coefficients of the time-domain equalizer being provided such that:\n(a) energy of an effective impulse response heff(n) of at least the communications channel combined with the time-domain equalizer is substantially concentrated in a second band of V+1 samples, whereby the second band of samples is shorter than the first band of samples; and\n(b) a variance in a frequency spectrum of the output additive noise signal of the time-domain equalizer is controlled.\nThe time-domain equalizer of the present invention preferably has an input additive noise signal, the T coefficients of the time-domain equalizer being provided as a function of both the input additive noise signal and an estimate of the impulse response h(n).\nThe receiver is preferably operable to compute the T coefficients of the time-domain equalizer as a function of balancing (i) a degree to which the energy of the effective impulse response heff(n) is concentrated: in the second band of V+1 samples; and (ii) a degree to which the variance in the frequency spectrum of the output additive noise of the time-domain equalizer is reduced.\nThe receiver is preferably operable to vary the degree to which the energy of the effective impulse response heff(n) is concentrated in the second band of V+1 samples; and the degree to which the variance in the frequency spectrum of the output additive noise of the time-domain equalizer is reduced.\nThe receiver is preferably operable to compute the T coefficients of the time-domain equalizer by evaluating the following system of equations:\nxe2x80x83min¦wwT(HoutTHout+xcex2outR)w, subject to wT(HinTHin+xcex2inR)w=1, where\nw is a Txc3x971 matrix representing the T coefficients of the time-domain equalizer,\nHout is an (M+Txe2x88x92Vxe2x88x922)xc3x97T matrix representing samples of the estimated impulse response h(n) of the communications channel which produce M+Txe2x88x92Vxe2x88x922 samples of the effective impulse response heff(n) not containing concentrated energy when matrix w is multiplied by Hout,\nHin is a (V+1)xc3x97T matrix representing samples of the estimated impulse response h(n) of the communications channel which produce the V+1 samples of the effective impulse response heff(n) having concentrated energy when matrix w is multiplied by Hin,\nR is a Txc3x97T additive noise correlation matrix constructed from the input additive noise signal, and\nxcex2out and xcex2in are scalars which vary the degree to which the energy of the effective impulse response heff(n) is concentrated in the second band of V+1 samples; and the degree to which the variance in the frequency spectrum of the output additive noise of the time-domain equalizer is reduced.\nThe receiver is preferably operable to obtain w by solving w=(GT)xe2x88x921emin, where G=HinTHin+xcex2inR, and emin is an eigenvector corresponding to the smallest eigenvalue of (G)xe2x88x921(HoutTHout+xcex2outR)(GT)xe2x88x921.\nThe present invention also provides a method of determining T coefficients of a finite impulse response time-domain equalizer having input and output additive noise signals, the time-domain equalizer being employed in a receiver for receiving a discrete multi-tone signal over a communications channel having an impulse response h(n), energy of the impulse response being substantially concentrated in a first band of samples, the method comprising the step of:\ncomputing the T coefficients of the time-domain equalizer such that:\n(a) energy of an effective impulse response heff(n) of at least the communications channel combined with the time-domain equalizer is substantially concentrated in a second band of V+1 samples, whereby the second band of samples is shorter than the first band of samples; and\n(b) a variance in a frequency spectrum of the output additive noise signal of the time-domain equalizer is controlled.\nThe method of the present invention also preferably comprises the step of balancing (i) a degree to which the energy of the effective impulse response heff(n) is concentrated in the second band of V+1 samples; and (ii) a degree to which the variance in the frequency spectrum of the output additive noise of the time-domain equalizer is reduced.\nThe method of the present invention also preferably comprises the step of varying the degree to which the energy of the effective impulse response heff(n) is concentrated in the second band of V+1 samples; and the degree to which the variance in the frequency spectrum of the output additive noise of the time-domain equalizer is reduced.\nIt is preferred that the method of the present invention also include the steps of:\nestimating the impulse response h(n) of the communication channel;\nestimating the input additive noise signal to the time-domain equalizer;\nproviding a Txc3x971 matrix w representing the T coefficients of the time-domain equalizer;\nproviding a (V+1)xc3x97T matrix Hin representing samples of the estimated impulse response h(n) of the communications channel which produce the V+1 samples of the effective impulse response heff(n) having concentrated, energy when matrix w is multiplied by Hin; and\nproviding an (M+Txe2x88x92Vxe2x88x922)xc3x97T matrix Hout representing samples of the estimated impulse response h(n) of the communications channel which produce M+Txe2x88x92Vxe2x88x922 samples of the effective impulse response heff(n) not containing concentrated energy when matrix w is multiplied by Hout.\nIt is also preferred that the method of the present invention include the steps of:\nproviding a Txc3x97T additive noise correlation matrix R from the input additive noise signal, and\ndetermining scalars xcex2out and xcex2in which control a degree to which the energy of the effective impulse response heff(n) is concentrated in the second band of V+1 samples; and a degree to which the variance in the frequency spectrum of the output additive noise of the time-domain equalizer is reduced\nIt is also preferred that the step of computing the T coefficients of the time-domain equalizer is obtained by solving the following system of equations:\nmin¦wwT(HoutTHout+xcex2outR)w, subject to wT(HinTHin+xcex2inR)w=1.\nIn the method of the present invention, it is also preferred that w is obtained by solving w=(GT)xe2x88x921emin, where G=HinTHin+xcex2inR and emin is an eigenvector corresponding to the smallest eigenvalue of: (G)xe2x88x921(HoutTHout+xcex2outR) (GT)xe2x88x921.\nOther features and advantages of the present invention will become apparent from the following description of the invention which refers to the accompanying drawing."}
{"text": "The majority of cars and other vehicles contain in-vehicle audio systems containing relatively high quality audio components such as amplifiers and speakers. Such in-vehicle audio systems typically include radios, and cassette, CD and minidisk players, and the like. However, the range of audio content available to a listener of in-vehicle audio devices is limited.\nA listener to a radio is limited temporally to the programming available at that time. If a listener misses a desired radio program at its normal broadcast time, he is unable to listen to that program later, at a time convenient for him, unless he recorded it during broadcast, A radio listener is also limited geographically to the programming broadcast from a within-range radio transmitter. It is currently not possible for a radio listener to travel abroad and still receive broadcasts from his home country's local radio stations on an in-vehicle radio. There is thus a need for a method for listening to a radio program via an in-vehicle radio system, while the radio program is not being currently transmitted on an in-range transmitter, without the listener having to pre-record the program.\nA listener to in-vehicle audio devices, such as a cassette or CD player, is limited to the selection of prerecorded audio content, typically music, that he has available in the appropriate format.\nA large amount of audio content, including both music and spoken content, is available over the Internet. The majority of in-vehicle audio systems are not Internet-enabled, however.\nA large number of listeners to in-vehicle audio devices also possess cellular phones. New generations of cellular phones are Internet-enabled, allowing audio content to be downloaded from the Internet. The sound quality from cellular phones is typically poor as they are designed to minimize size and weight and are intended to be used by a single user, while being held close to the ear. Cellular phones are not designed for broadcasting music to a number of listeners. There is thus also a need for apparatus and a method for improving the sound quality of audio content downloaded from a network such as the Internet, through a cellular phone."}
{"text": "A variety of methods exist for casting materials of various sorts onto a moving substrate. In the photographic film manufacture process it is conventional practice to cast, for instance, cellulose acetate films by a solvent cast method. In the solvent cast method a solution of cellulose acetate in a solvent is poured onto a moving casting surface. After sufficient evaporation of the solvent, the cellulose acetate film is peeled from the casting surface. A near infra-red (NIR) gauge is used for measuring the solvent content of the acetate film while on the casting surface. For measurement purposes, the NIR gauge needs to be mounted about 0.250 inches to about 0.375 inches away from the casting surface. The casting surface has a mirror finish and can easily be damaged. Protuberances or obstructions on film exceeding the dimensions of the spacing between the gauge and casting surface are caught between the NIR gauge and the casting surface. This generally results in damage to the finish of the casting surface and to the probe.\nWhile there are no known attempts to solve the specific problem confronted by Applicants, U.S. Pat. No. 4,865,872, titled \"Strip Inspecting Apparatus and Associated Method\" by Pellatiro, Sep. 12, 1989, discloses a strip inspection apparatus wherein a linear sensor array is positioned above the surface of a strip for determining the position of the strip edges. Further, U.S. Pat. No. 4,797,301, titled \"Method of and Apparatus for Applying a Fluid to a Bulk Commodity\" by Ardley et al., Jan. 10, 1989 discloses a combination pivotally mounted apparatus for applying a fluid and for monitoring the thickness thereof to a material being conveyed. An obvious shortcoming of each of the aforementioned devices for solving Applicants' problem is that neither provides for a flexural assemblage that yields to contact by a surface protuberance thereby avoiding damage to the surface undergoing inspection.\nTherefore, a need persists for a flexural probe assemblage for examining a sensitive web surface that can be arranged very close to the sensitive surface and can then pivot away from the sensitive surface when engaged by a protuberance on the surface without causing damage to the surface."}
{"text": "The present disclosure relates to a projector. More specifically, the present disclosure relates to a presentation projector that dynamically masks out one or more objects positioned between the projector and a screen.\nDuring a typical office presentation, a digital data file such as a slideshow may be loaded onto a computing device and displayed on a screen via a projector. FIG. 1 illustrates an exemplary projector 100 projecting a presentation onto a screen 102. The projector 100 receives a video-in signal 104. The video-in signal 104 is processed by a liquid crystal display (LCD) panel controller 106. In a typical LCD projector, three LCD panels 108 are used, one each for projecting the colors red, blue and green. By opening or closing each pixel on the LCD panels 108, a variety of colors may be projected for each pixel by combining the outputs of the various LCD panels 108. Based upon the video-in signal 104, the LCD panel controller 106 determines the state of each pixel (i.e., opened or closed) for every pixel in the LCD panels 108. The LCD panel controller 106 transmits a control signal to the LCD panels 108 indicating the state (i.e., open or closed) for each pixel in the panel. A lamp 110 projects a strong beam of light through each of the LCD panels 108. The lamp 110 may be a high intensity lamp such as a metal halide lamp or other similar lamp. After the light emitted by the lamp 110 passes through the LCD panels 108, the light passes through one or more lenses 112, and the projected light 114 is directed by the lenses toward the screen 102.\nThe projector 100 may be a portable device placed on a desk or table at a low enough level that the light emitted from the projector may be projected at eye level of a person 116 walking in front of the projector. The person 116 walking in front of the projector may experience discomfort due to bright light shining in their eyes. Audience members may be distracted by the images projected onto the person 116 walking in front of the projector, as well as the shadow 117 cast by the person on the screen 102."}
{"text": "This invention relates to fluorescent microscopy, and more particularly to the rapid selection of filters for measuring fluorescence at different wavelengths or Stokes shifts.\nThe phenomenon of fluorescence describes light emission that continues only during the absorption of the excitation light by a chromophore or other conjugated molecule, which is capable of emitting secondary fluorescence. Reflected light fluorescence microscopy is a favored technique in fluorescence microscopy. This mode of fluorescence microscopy is also known as incident light fluorescence, fluorescence, or scopic fluorescence. The universal reflected light vertical illuminator is interposed between the observation viewing tubes and the nosepiece carrying the objectives. A key feature of fluorescence microscopy is its ability to detect fluorescent objects that are sometimes faintly visible or even very bright relative to the dark (often black) background. In order to optimize this feature, image brightness and resolution must be maximized. Multiphoton fluorescence microscopy is a powerful research tool that combines the advanced optical techniques of laser scanning microscopy with long wavelength multiphoton fluorescence excitation to capture high-resolution, three-dimensional images of specimens tagged with highly specific fluorophores.\nWhile fluorescence is described, the invention can be used with any type of photoluminescence, provided that the material under observation exhibits a Stokes shift.\nFluorescence microscopy can be used for specimens, which exhibit autofluorescence, and for specimens enhanced by fluorochromes, also called fluorophores. The basic task of the fluorescence microscope is to permit excitation light to irradiate the specimen and then to separate the much weaker re-radiating fluorescent light from the brighter excitation light. Thus, only the emission light reaches the eye or other detector. The resulting fluorescing areas shine against a dark background with sufficient contrast to permit detection. The darker the background of the non-fluorescing material, the more efficient the instrument. Therefore efficient filters are desired.\nWhen a fluorescing sample is observed with a fluorescence microscope. Ultraviolet (UV) light of a specific wavelength or set of wavelengths is produced by passing light from a UV-emitting source through the exciter filter. The filtered UV light illuminates the specimen, in this case a crystal of fluorspar, which emits fluorescent light when illuminated with ultraviolet light. Visible light emitted from the specimen, red in this case, is then filtered through a barrier filter that does not allow reflected UV light to pass. It should be noted that this is the only mode of microscopy in which the specimen, subsequent to excitation, gives off its own light. The emitted light re-radiates spherically in all directions, regardless of the direction of the exciting light.\nFor a given fluorochrome, the manufacturer indicates the wavelength for the peak of excitation/fluorescence intensity and the wavelength for the peak of emission/fluorescence intensity. To determine the emission spectrum of a given fluorochrome, the dye absorption maximum wavelength is found and the fluorochrome is excited at that maximum. The dye absorption maximum wavelength is usually the same as the excitation maximum. The absorption spectrum of a typical fluorochrome occurs where the relative intensity of absorption is plotted against relative wavelength. A monochromator is then used to scan the fluorescence emission intensity at successive emission wavelengths. The relative intensity of the fluorescence is measured at the various wavelengths to plot the emission spectrum. The excitation spectrum of a given fluorochrome is determined in a similar manner. The emission maximum is chosen and only emission light at that emission wavelength is allowed to pass to the detector. Then excitation is induced at various excitation wavelengths and the intensity of the emitted fluorescence is measured. There is usually an overlap at the higher wavelength end of the excitation spectrum and the lower wavelength end of the emission spectrum. This overlap of excitation and emission intensities/wavelengths must be eliminated in fluorescence microscopy. Elimination of the overlap of excitation and emission intensities/wavelengths is accomplished by means of appropriate selection of excitation filter, dichromatic beamsplitter, and excitation or barrier filter. Otherwise, the much brighter excitation light overwhelms the weaker emitted light and significantly diminishes specimen contrast.\nThe separation of excitation and emission wavelengths is achieved by the proper selection of filters to block or pass specific wavelengths. The design of fluorescence illuminators is based on control of excitation light and emission light by readily changeable filter insertions in the optical pathway on the way toward the specimen and then emanating from the specimen.\nA dichroic material is one which absorbs light polarized in one direction more slowly than light polarized at right angles to that direction. Dichroic materials are to be distinguished from birefringent materials, which may have different refractive indexes for two electric vectors vibrating at right angles to each other but similar, if negligible, absorptions coefficients. The term dichroic is also used to denote the change in color of a dye solution with change in concentration, to denote a color filter that has two transmission bands in very different portions of the visible region and hence changes color when the spectral distribution of the illuminating source is changed, and to denote an interference filter that appears to be of a different color when viewed in reflected or transmitted light.\nAs used in this disclosure, xe2x80x9cdichroic beamspiitterxe2x80x9d refers to a mirror, which reflects most light below a predetermined wavelength and transmits light above that wavelength. Thus, in one sense the invention may use a dichroic filter, which is an interference filter that appears to be of a different color when viewed in reflected or transmitted light.\nDichroic filters are filters having or showing two colors. This can be obtained with doubly-refracting crystals that exhibit different colors when viewed in different directions or with solutions that show essentially different colors in different degrees of concentration.\nIn addition, while dichroic materials are described, the invention may be used with birefringent materials and other types of filters.\nThe functions of a dichroic element in a system used to analyze fluorescent labels in samples include directing light from a light source toward the sample under observation, and transmitting light reflected from the sample. By limiting the reflection to a predetermined color, it is possible to specifically detect reflected light within a predetermined waveband.\nWhen using two or more fluorescent labels or fluorophores, it is necessary to provide dichroic beamsplitter mirrors for each of the fluorescent labels. It is further desired to detect these fluorescent labels separately. While the fluorescent labels are detected separately, there should not be a significant time difference between the detection of the different fluorescent labels. Optimally, it would be desirable to simultaneously detect the different fluorescent labels. If simultaneous detection is not practical, then the sequence of detecting the different fluorescent labels should be rapid enough that the end effect is similar to that of simultaneous detection.\nIn a particular type of fluorescence microscope, a dichroic beamsplitter mirror is used in combination with an excitation filter and an emission filter, sometimes referred to as a barrier filter. The excitation filter filters the excitation light to provide a monochromatic output, and the emission filter filters light transmitted through the dichroic mirror. The dichroic beamsplitter mirror and the emission filter are selected according to the Stokes shift of the fluorophores under observation.\nCurrent fluorescent microscope designs employ an incident light or fluorescent design where a dichroic beamsplitter or chromatic mirror is mounted in a filter cube at a 45 degree angle to the excitation light path. The beamsplitter or chromatic mirror is used to reflect shorter excitation wavelengths of light onto the specimen while passing longer emission wavelengths to the eyepieces or camera. The beamsplitter has a cutoff wavelength or frequency which is between an excitation wavelength or frequency and an emission wavelength or frequency of the fluorescent material being sensed.\nIn addition to the dichroic beamsplitter, a monochromatic light source is provided at an excitation frequency, and a notch filter at an emission frequency is used at the sensing unit. The wavelengths of the monochromatic light source and the filter at the sensing unit generally correspond to the excitation wavelength and the emission wavelength of the fluorescent material being sensed. In the typical implementation the monochromatic light source is provided by using a notch filter to block emissions from a high intensity lamp at wavelengths above and below a selected excitation wavelength.\nMany fluorescence microscopy applications require two or more fluorescent labels to be present in the specimen. Each label has its own excitation and emission spectra, and thus requires different excitation and emission filters, as well as a different dichroic beamsplitter.\nTo date two approaches to using multiple fluorophores have been employed:\n1) Use dichroic beamsplitters that have multiple cut-off wavelengths. Thus a single dichroic can be used with multiple fluorophores. However because of bandwidth restrictions, total light throughput is reduced, thus creating longer exposure times when working with a camera. Longer exposure times translate into longer acquisition duty cycles, which is problematic in paradigms requiring repetitive high-speed data acquisition, such as screening applications and applications using living cells.\n2) Motorized filter cube changers. Several commercially available microscopes employ motorized filter cube changers. These allow the use of a single dichroic for each fluorophore. However switching time is slow (12 seconds) which creates problems when using multiple fluorophores in paradigms requiring repetitive high-speed data acquisition. The slow switching time is the result of the significant mass of the motorized filter cubes. Increased speed of switching would require a substantial increase in power and control of reaction forces caused by the movement of the cubes during switching.\nDichroic beamsplitter mirrors are available with two or more cutoff wavelengths. This results in an ability to use such a beamsplitter with a first set of excitation and emission filters corresponding to a first cutoff wavelength to sense a first fluorophore, and to use the same beamsplitter mirror with a second set of excitation and emission filters with a second combined cutoff wavelength to sense a second fluorophore. This becomes difficult to accomplish if the beamsplitter mirror is provided in a filter cube.\nAccording to the present invention, an optical system for a microscope, which includes an illumination source, microscope optics for obtaining a magnified image, and a beamsplitter changer is able to rapidly effect changes in Stokes shift wavelengths being detected. A plurality of dichroic beamsplitter elements (xe2x80x9cdichroic beamsplittersxe2x80x9d) are arranged about a center axis at a preferred angle such as 45xc2x0 to the lightpath of the microscope optics, or optical pathway. The beamsplitter elements are rotable about the center axis to allow positioning a selected one of the dichroic beamsplitter elements in a part of the light path between the illumination source and the microscope optics, with the light path resulting in reflection of light from the illumination source to a specimen under observation, and transmitting light received through the microscope optics from the specimen. When the selected dichroic element corresponds to a fluorescent material under observation, the dichroic beamsplitter reflects excitation light at an excitation wavelength corresponding to the fluorescent material under observation and the dichroic element transmits light shifted in wavelength by a Stokes shift of the fluorescent material under observation. Rapid positioning of the dichroic filter is achieved by a motor and motor controller, which rotates the filter about the center axis."}
{"text": "The use of unsaturated polyester resins in the molding of reinforced products enjoys broad application in the manufacture of automotive, industrial, and home products. With the advent of gasoline shortages, the automotive industry has increased its efforts to replace metal parts with fiber-reinforced composites to reduce weight and gasoline consumption. In this respect, it has been an important objective to improve physical properties of the polyester composites.\nTo improve crack resistance of fiber-reinforced composites, polyester resins of the type used in SMC have been modified by addition of reactive liquid polymers, such as vinyl terminated butadiene-acrylonitrile. The fracture surface energy of the modified polyester resin was thus improved by better than 9 times at a loading of 10 parts of the liquid polymer per 100 parts of about 40/60 alkyd/styrene resin and better than 3 times for a 42/58 alkyd/styrene resin. Although the toughening effect of the liquid polymers was very attractive, the resulting modified polyester resins presented mixing problems due to the incompatiblity of the liquid polymers with the polyester resins.\nNumerous attempts have been made to improve mechanical properties of filled polyester composites by admixing fillers, polyester resins, reactive liquid polymers, and other additives with the polyester resins, however, such attempts have been only marginally successful or totally unsuccessful due to inherent brittleness and incompatibility of the materials, or other problems.\nU.S. Pat. No. 3,518,221 to Kenyon, et al., discloses a molding composition and a process for preparing same by coating reinforcing particles with a first thermosettable resin, and then embedding the coated particles in a second thermosettable matrix resin, the second matrix resin being more rigid than the first resin. The first and second resins can be same or different as long as the second resin is more rigid than the first resin. The resins can be selected from epoxy, polyester, phenolic, amide, imide, amine, or urethane resins. In a preferred embodiment, the resins are selected from epoxide resins such as the Epons and Araldites. A particularly suitable commercially available resin is Epon-815 having average molecular weight of about 330. Such compositions can be molded to form shaped articles having very good tensile strength and modulus.\nU.S. Pat. No. 3,956,230 to Gaylord discloses preparation of a composition by coating hydroxyl-containing filler particles with a first thermoplastic polymer containing labile atoms in the presence of an ethylenically unsaturated carboxylic acid or anhydride coupling agent under conditions which will generate free radicals on the thermoplastic polymer. Free radicals are formed by using a catalyst, such as a peroxide. The acid is thereby coupled onto the thermoplastic polymer and is reacted with the hydroxyl groups on the filler particles by means of esterification and hydrogen bonding. The coated filler particles are then dispersed in a thermoplastic polymer matrix, which can be same or different as the first thermoplastic coating resin. Suitable thermoplastic resins, which are disclosed at bottom of col. 2 and in col. 3 of the Gaylord patent, include polyurethanes which may be obtained by reaction of a diisocyanate with a polyol such as polyethylene oxide, polypropylene oxide, polytetramethylene glycol, hydroxyl-terminated polyesters, hydroxyl-terminated polyisobutylene and hydroxyl-terminated polybutadiene.\nAs the Kenyon, et al. and Gaylord patents will verify, which patents are discussed above, the prior art is replete with disclosures of epoxy resins used as a coating on particulate filler particles which are thereafter mixed with a matrix resin to form molding and other compositions. Epoxy resins of lowest molecular weight, such as Epon.RTM. 815 and 828, are based on diglycidyl ether and bisphenol A and have number average molecular weight (Mn) of about 380 to 400 and heat distortion temperature (HDT) of about 65.degree. to 75.degree. C., see p. 66 of the text entitled \"Epoxy Resins\" by Lee and Neville, published by McGraw-Hill Book (1957). Higher molecular weight resins have higher HDTs. HDT depends to a degree on the curing agent used. For instance, diethylene tetramine curing agent will yield epoxy resins having low HDT whereas aromatic amines as curing agent will give epoxy resins having high HDT of about 150.degree. C. Generally, HDT and Tg of a material will be within 10.degree. C. of each other, with HDT being lower than Tg. Comparison of the hereinclaimed functionally terminated reactive liquid polymers with the lowest molecular weight epoxy reins shows that Tgs of the epoxy resins are at least 50.degree. C. higher than Tgs of the reactive liquid polymers. Of course, for higher molecular weight epoxy reins, the difference will be greater."}
{"text": "Field of the Invention\nThis invention relates to production equipment for oil wells and more particularly, to sucker rod guides for mounting in spaced relationship on the sucker rod string of a pumping well. The sucker rod guides of this invention are each characterized by a generally cylindrically-shaped, smooth guide body having tapered top and bottom portions and a longitudinal, grooved bore which is preferably slightly undersized with respect to the sucker rod to which the sucker rod guides are attached. A tapered, inwardly expanding, longitudinal body slot extends through the entire guide body of the sucker rod guide and communicates with the longitudinal bore, in order to provide a means for mounting one or more sucker rod guides on a sucker rod at a specific location. In a preferred embodiment, the sucker rod guide of this invention is constructed by extruding a molten, high density, self-lubricating polyethylene plastic through a die to define a slotted cylinder having a smooth bore and a tapered longitudinal slot. Grooves may then be added to the bore by using a reaming tool, to facilitate a more firm attachment to the sucker rod. Alternatively, the extruded, slotted cylinder can be centrally drilled to provide the longitudinal bore and the spaced grooves.\nIt has surprisingly been found that extruded molten polyethylene creates an ultra-high density, round bar stock which produces a sucker rod guide that yields an abrasion factor which is far superior to that of conventional sucker rod guides. In fact, it has been found by analysis and testing that extruded high density polyethylene plastic contains approximately two million molecules per square inch of extruded bar stock material, a density factor which accounts for the superior abrasion and self-lubricating characteristics of the extruded plastic material. It has also been found that the polyethylene sucker rod guides of this invention are not adversely affected by corrosive hydrogen sulfide, salt water and other fluids and compounds normally found in an oil well.\nSucker rod guides of various description and composition are designed to fit on sucker rods used to pump oil wells, in order to eliminate, or at least greatly reduce, many of the down-hole problems which are characteristic of production equipment in oil wells. These guides are generally characterized by a coefficient of friction when wet which is lower than that of metal and they operate to increase the overall pumping efficiency of the wells, while at the same time prevent undesirable metal-to-metal contact between the reciprocating or spinning rods and the stationary tubing. Wear on sucker rod couplings used to make up the down-hole sucker rod string is also minimized, thus reducing the required inventory of costly rod parts. Tubing wear, often unseen until failure occurs, is also reduced, because the sucker rod guides receive the wear rather than the expensive tubing. The sucker rod guides also function as bearings to centralize the sucker rods in the tubing and even when the tubing buckles in a well during upstroke of the pumping unit, the sucker rods cannot contact the tubing due to the spaced positioning of the sucker rod guides. Polish rod loads are also reduced because of the lower friction and less abrupt \"stress reverse\" which is realized when using sucker rod guides. Accordingly, a properly designed sucker rod installation can result in significant savings in both equipment replacement and service costs in a pumping oil well. Fewer \"pulling jobs\", greater pumping efficiency and wells that stay in the pumping mode for longer periods of time, are proven results obtained from the use of sucker rod guides.\nSucker rod guides of various design, size and materials of construction are well known in the art. Typical materials of construction are neoprene rubber and nylon and these materials are commonly used with metal inserts which encase and line the sucker rod bore, in order to better grip the sucker rod and maintain the sucker rod guide in a selected position on the sucker rod. However, it has been found that friction between neoprene rubber sucker rod guides and the tubing as the sucker rod guide reciprocates with the sucker rod string inside the tubing, sometimes generates heat which may result in a fairly rapid deterioration of the neoprene material, thereby necessitating frequent \"pulling jobs\" in order to replace the guides. Furthermore, it has been found that nylon sucker rod guides are brittle and are sometimes difficult to mount on a sucker rod without breaking, especially in cold weather. Sucker rod guides are typically secured to the respective lengths of sucker rod in spaced relationship, in order to space the sucker rod string from the tubing and protect both the sucker rod, sucker rod couplings and the tubing from excessive wear during the pumping operation. Since the reciprocating travel of each sucker rod and sucker rod guide may be from approximately 3 feet to about 20 feet or more and this travel occurs at a rate of about 15 strokes per minute on the average, the sucker rod guide material should be self-lubricating or easily lubricated by the well fluid and must have good wear characteristics, in order to minimize the frequency of maintenance.\nAccordingly, it is an object of this invention to provide a new and improved plastic sucker rod guide for mounting on the sucker rod of an oil well and preventing, or at least minimizing, metal-to-metal contact between the sucker rod and the tubing.\nAnother object of this invention is to provide a new and improved molded polyethylene sucker rod guide which is characterized by an exceptionally high molecular density and good self-lubricating characteristics and is designed to tightly seat on the sucker rod of a pumping well at a specific location, with minimum displacement from the point of installation, to space the sucker rod from the tubing.\nStill another object of this invention is to provide a sucker rod guide for mounting on the sucker rod of a pumping unit in an oil well, which sucker rod guide includes a generally cylindrically-shaped, extruded polyethylene body having a longitudinal, smooth or grooved sucker rod bore which is slightly undersized with respect to the sucker rod to which it is attached and further including an inwardly expanding, tapered body slot extending longitudinally through the sucker rod guide body and communicating with the sucker rod bore, to facilitate mounting the sucker rod guide on the sucker rod.\nA still further object of this invention is to provide a generally cylindrically-shaped, extruded, high density polyethylene sucker rod guide provided with a longitudinal, circumferentially grooved sucker rod bore which is slightly undersized with respect to the sucker rod to which the guide is attached and having a longitudinal, tapered slot extending through the sucker rod guide body and communicating with the sucker rod bore, for mounting the sucker rod guide on the sucker rod."}
{"text": "Cycle exercise machines, sometimes also known as \"stationary bicycles\", have been known for some time. These exercise machines typically include a seat much like a bicycle seat, for supporting a person sitting upright on the machine, and a foot-powered crank arrangement mounted below the seat. This crank arrangement in prior-art cycle exercisers includes a pair of crank arms extending outwardly from a rotatable shaft, with foot pedals mounted at the free ends of the crank arms. The crank arrangement is mounted for receiving the feet of a person seated on the exercise machine. A drive chain or belt transfers the rotary motion of the crank arrangement to an energy-absorbing device such as a flywheel or rotor, which provides a load force against which the exerciser expends energy while pedaling the crank mechanism. One such cycle exerciser is disclosed in U.S. Pat. No. 4,188,030 to Hooper.\nBecause cycle exercisers of the prior art use a crank mechanism which the user pedals while doing work on the machine, these exercise machines cannot provide a constant transfer of power from the user to the load throughout the entire stroke of each pedal. Crank pedal mechanisms have a null at the top and bottom of each stroke, where the effective length of the pedal crank arms diminishes to zero length as the crank mechanism passes through the dead center position. The effective lever arm of the pedals then increases sinusoidally to the maximum effective length as the pedal crank arms pass through the 90.degree. position, i.e., half-way between the top to the bottom of each stroke. Because this effective lever arm constantly changes, the rate at which the exerciser can effectively expend energy doing work on the load device likewise changes throughout each stroke of the pedal. This variation lessens the possible maximum efficiency of cycle exercisers, as the rate of energy transfer for each stroke is maximized only momentarily during each stroke.\nAnother disadvantage of existing cycle exercisers, and in particular the crank mechanism used with those exercisers, pertains to the length of stroke. The stroke length is determined by the length of the crank arms used in the crank mechanism, and this length is not readily changeable. Changing the effective stroke length requires either replacing the entire crank assembly, or providing crank arms having variable positions for attaching the pedals on the crank arms. In either case, these alternatives are expensive and require the service of a technician to modify the length of the crank arms for users having significantly different lengths of stroke, namely, shorter or longer legs, or those of substantially different athletic ability."}
{"text": "1. Field of the Invention\nThe present invention relates to a disc cartridge which rotatably accommodates a disc-shaped disc medium used as a recording/playback medium of an information processing device or the like, and which prevents entry of dust or the like.\n2. Description of the Related Art\nA disc-shaped disc medium, such as an optical disc, a magneto-optical disc, or the like, is used, for example, as a portable recording/playback medium for a computer. When recording or playback of the disc medium is carried out, laser beam is illuminated onto the recording surface while the disc medium is rotated in a state of being loaded in a drive device. In this way, recording of information is carried out by pit formation or phase changes or magnetization or the like due to decomposition of the dye layers of the recording surface, or recorded information is played back on the basis of differences in the reflectivity or the polarizing angle of the laser beam.\nFurther, in such disc media, shortening of the wavelength of the laser beam which is illuminated onto the recording surface has been proposed in order to increase the recording capacity. When recording or playback of information is carried out by using this short-wavelength laser beam (e.g., a blue-violet color laser), in order to suppress the attenuation of the laser beam due to the cover layer which protects the recording surface of the disc medium, the cover layer must be made thin. By making the cover layer thin in this way, the diameter of the laser beam at the surface of the cover layer (the surface which is exposed to the exterior) decreases, and it is not possible to ignore the effects of dust and the like adhering to the surface of the cover layer.\nThus, disc cartridges, which accommodate a disc medium in a case and prevent adhesion of dust or the like onto the disc medium, are employed. Such a disc cartridge is formed to include an opening, which is for exposing to the exterior a center hole portion (center core portion) provided at the central portion of the disc medium and a portion of the recording surface (the cover layer), and a shutter member which opens and closes this opening.\nIn this way, in the disc cartridge, generally, by closing the opening by the shutter member, entry of dust or the like into the disc cartridge, i.e., adhesion of dust or the like onto the disc medium, is prevented, and by opening the opening as the disc cartridge is loaded into the drive device, the center hole portion (the center core portion) can be held by a rotation spindle shaft, and the laser head can approach the recording surface (laser beam can be illuminated).\nFurther, in such a disc cartridge, in the aforementioned state in which the opening is open, there is the need to ensure a space for accommodating the shutter member. Also, it is desired to make the disc cartridge compact and thinner. Thus, to aim for compactness and thinness, a disc cartridge in which the shutter member is divided in two has been disclosed in Japanese Patent Application Laid-Open (JP-A) No. 2000-30394 and the like.\nIn the disc cartridge of JP-A 2000-30394, the entire opening is closed in a state in which the end surfaces of a first shutter member, which mainly closes the opening portion for laser head insertion (approach), and a second shutter member, which mainly closes the opening portion for insertion of the rotation spindle shaft, abut one another. In a opening state of the opening, the first shutter member and the second shutter member rotate in different directions while above mentioned abut surfaces are being apart each other. The first shutter member and the second shutter member are accommodated within accommodating spaces which substantially correspond to the projection surface area of the disc-shaped disc medium. Thus, this disc cartridge aims for a reduction in size and a reduction in thickness, while preventing adhesion of dust and the like onto the disc medium.\nAt one of wall surfaces of this disc cartridge, a guide groove for guiding at a time of inserting the disc cartridge into the drive device is provided in parallel to the inserting direction. A lock releasing lever of a lock mechanism which regulates (stops) movement of the shutter member, and an activating member which moves the shutter member to an opening direction protrude in the guide groove. A convex portion, which engages with the lock releasing lever and the activating member, protrudes from the drive device.\nAccordingly, when such disc cartridge is inserted into the drive device, the convex portion enters in the guide groove in accompany with the insertion of the disc cartridge into the drive device. First, the convex portion abuts the lock releasing lever and a lock state of the shutter member is released. Then, when the disc cartridge is inserted further into the drive device, the convex portion abuts the activating member and the shutter member is moved in the opening direction. In this way, the opening is opened and access to the disc media becomes possible.\nIn this way mentioned above, the shutter member which opens and closes the opening of the disc cartridge is moved by operating the lock releasing lever and the activating member, which are disposed at the one side wall of the disc cartridge. However, in this case, because of means which can move the shutter member so as to open the opening being provided at one side wall, there is a problem that countermeasure against the shutter member being inadvertently opened does not suffice.\nThe object of the present invention is to provide a disc cartridge which prevents a shutter portion from being opened and closed inadvertently.\nIn order to achieve the above object, the disc cartridge of a first aspect of the present invention is a disc cartridge having a shutter portion which can open and close an opening for access to a disc-shaped disc medium, the opening being formed in at least one surface of a case which rotatably accommodates the disc medium, wherein two groove portions are provided at wall surfaces of the case, and an operating portion of a lock mechanism for the shutter potion is provided in one of the groove portions, and an activating portion for opening and closing operation of the shutter portion is provided in another of the groove portions.\nThe disc cartridge of a second aspect of the present invention is a disc cartridge having a shutter portion which can open and close an opening for access to a disc-shaped disc medium, the opening being formed in at least one surface of a case which rotatably accommodates the disc medium, the disc cartridge comprising: a first groove portion which is formed on one of wall surfaces of the case; a second groove portion which is formed on another of the wall surfaces of the case; a lock portion which regulates movement of the shutter portion; and an activating portion for moving the shutter portion, wherein the activating portion protrudes from the first groove, and an operating portion of the lock portion protrudes from the second groove.\nIn the disc cartridge of a third aspect of the present invention according to the second aspect, a protruding portion which protrudes from the second groove is formed at one end portion of the operating portion, and an engaging portion for engaging with the shutter portion is formed at another end portion of the operating portion.\nIn the disc cartridge of a fourth aspect of the present invention according to the third aspect, the shutter portion comprises a first shutter member and a second shutter member, and the activating portion is provided at an one end portion of the first shutter member and an engaged portion which is engaged with the engaging portion of the lock portion is provided at another end of the first shutter member.\nIn the disc cartridge of a fifth aspect of the present invention according to the second aspect, the first groove portion and the second groove portion are formed respectively at side walls, which oppose each other, of the case.\nIn order to open the opening intentionally, the operating portion of the lock mechanism and the activating portion for opening and closing operation of the shutter portion are operated at the same time. However, in the present invention, the operating portion and the activating portion are provided in the respective different groove portions, namely, operating of the shutter member is not facile. Accordingly, it can be prevent that the opening is opened inadvertently.\nIt is preferable that the case is formed by joining a top shell and a bottom shell, and the shutter portion is provided at an outer side of the case. Further, it is preferable that a covering plate portion, in which an opening corresponding to the above mentioned opening is formed, is provided, and the shutter portion is movably accommodated between the covering plate portion and the case.\nFurther, it is preferable that the lock mechanism (the lock portion) is structured such that a center portion of the operating portion of lock mechanism is rotatably supported, a protruding portion provided at one end of the operating portion protrudes in the one of the groove portions and an engaging portion provided at another end of the is engaged with a cut out portion formed at the shutter portion. In this way, the lock mechanism can be operated preferably. Also, because the cut out portion is a shutter window portion forming a part of the opening, due to utilizing this window portion, it becomes not necessary that another engaged portion for engaging with the engaging portion of the operating portion of the lock mechanism is provided at the shutter portion. Accordingly, the cost does not increase."}
{"text": "This invention relates to electronic data processing systems and more particularly to a central processing unit for use in such systems."}
{"text": "The present invention relates generally to toggle action hand tools, and more specifically to an adjustable length locking bar clamp with a locking pliers type, toggle action clamping mechanism.\nLocking pliers type hand tools with toggle action clamping mechanisms are generally known as locking pliers or vice clamps. These conventional tools usually include a fixed handle having a fixed jaw extending from one end of the fixed handle, a movable handle, and a toggling linkage arrangement interconnecting the fixed handle and the movable handle. The toggling arrangement typically includes a movable jaw extending from the toggling linkage arrangement in such a way that when the locking pliers are toggled from a first open position to a second toggled position the movable jaw is moved toward the fixed jaw.\nSince the original development of the locking pliers a wide variety of jaw arrangements have been developed for all types of applications. For example, U.S. Pat. No. 4,850,254 describes a locking hand tool including locking pliers in which a user adjustable jaw and a movable jaw are moved toward one another in a direction generally perpendicular to the handles of the locking pliers when the locking pliers are toggled into the locked or toggled position. This arrangement provides an adjustable clamp in which the locking pliers extend perpendicular to the direction of clamping action. Although this a useful arrangement in some situations, there are many situations in which the generally perpendicular location of the locking pliers interferes with the work being performed. The present invention discloses a toggle action quick release bar clamp having a specific jaw arrangement for toggle action locking pliers which provides a bar clamp with a clamping action direction that is along the longitudinal axis of the locking pliers handle or parallel with the handle."}
{"text": "This invention pertains to implantable medical devices such as cardiac pacemakers and implantable cardioverter/defibrillators. In particular, the invention relates to a system and method enabling an implantable medical device to communicate with and control an external drug delivery device.\nImplantable medical devices are commonplace today, particularly for treating cardiac dysfunction. Cardiac pacemakers, for example, are implantable medical devices that replace or supplement a heart\"\"s compromised ability to pace itself (i.e., bradycardia) due to chronotropic incompetence or a conduction system defect by delivering electrical pacing pulses to the heart. Implantable cardioverter/defibrillators (ICD\"\"s) are devices that deliver electrical energy to the heart in order to reverse excessively rapid heart rates (tachycardia) including life threatening cardiac arrhythmias such as ventricular fibrillation. Since some patients have conditions that necessitate pacing and also render them vulnerable to life-threatening arrhythmias, implantable cardiac devices have been developed that combine both functions in a single device.\nMost pacemakers today are operated in some sort of synchronous mode where the pacing pulses are delivered in a manner that is dependent upon the intrinsic depolarizations of the heart as sensed by the pacemaker. ICD\"\"s must also sense the electrical activity of the heart in order to detect an arrhythmia that will trigger delivery of the shock pulse in an attempt to reverse the condition. Such sensing information can be used to initiate another mode of therapy, and efforts have been made in the past to combine automatic drug delivery by an implantable drug delivery system with either pacemakers, ICD\"\"s, or both in order to treat cardiac arrhythmias.\nImplantable drug delivery systems suffer from a number of disadvantages, however, when compared with an external drug delivery device. Although the drug reservoir of an implantable delivery device can be replenished, it is difficult to change the drug once it is put into the reservoir, making patient management difficult in cases where a patient\"\"s condition either changes or otherwise requires a change of medication. In addition, drugs degrade over time. Finally, there is the risk of leakage from the implanted reservoir, the consequences of which can range from an annoyance to a medical emergency. For these reasons, drug delivery from an external device is preferred in many situations.\nIn order to control the delivery of drugs or other therapies by an external device, the implanted device must be capable of transmitting command and control information to the external device. This is the primary problem with which the present invention is concerned.\nThe present invention relates to a drug delivery system and method in which an implantable medical device communicates with and controls an external drug delivery device with low energy potential signals. The implantable device generates potential signals by operating a current source to cause corresponding electrical potentials that can be sensed at the skin surface by the external drug delivery device. The current source is operated so as to generate a carrier waveform that can be modulated with digitally encoded command information. The external drug delivery device includes electrodes at the skin surface for sensing potentials and circuitry for demodulating the sensed carrier waveform. The drug delivery device may then decode the digital data to extract command information therefrom and deliver a drug accordingly.\nIn accordance with the invention, the carrier waveform is digitally modulated with the digitally encoded information by varying the amplitude or frequency of the carrier waveform using, for example, amplitude shift-keying or frequency shift-keying. In a particular embodiment, a digital pulse train is modulated with the digitally encoded information by varying the frequency, width, or position of the pulses. The pulse train is then used to amplitude modulate the carrier waveform.\nCertain implantable medical devices, such as rate-adaptive pacemakers, may use an impedance technique for measuring minute ventilation and/or cardiac stroke volume. In that technique, an oscillating current is made to flow between two electrodes located within the thorax, and the impedance between the electrodes is measured. In accordance with the invention, the impedance measuring current may be used as the carrier waveform and modulated with digitally encoded information by the implantable device for transmission to the external drug delivery device.\nThe present invention may be incorporated into a system where the implantable device is a cardiac device such as an implantable cardioverter/defibrillator, cardiac pacemaker, or combination device. A dose of a drug is then delivered by the drug delivery device to a patient upon detection of a particular medical condition such as an arrhythmia. In one embodiment, the drug delivery device is an electrically modulated transdermal drug delivery device designed for external affixation to a patient\"\"s skin surface at a suitable location. The electrically modulated transdermal delivery injector may have one electrode with a drug reservoir in contact with the skin, another electrode also contacting the skin, and a voltage source for imposing a voltage between the electrodes. The drug delivery device also has a data communications interface for receiving command signals from the implantable cardiac device, and circuitry for controlling the delivery of the drug in accordance with the command signal."}
{"text": "The invention relates to filters, including air filters, and more particularly to a low cost mounting and retention system.\nAn air filter traditionally includes three components, namely: (1) a filter element; (2) a housing base; and (3) a housing cover mounted to the base and enclosing and/or sealing and/or retaining the filter element in the housing. The parent invention of the noted '144 application eliminates the need for the third component, and makes the use of a cover optional. In the preferred embodiment, the noted parent invention eliminates the need for the third component by directly attaching/retaining and sealing the first component to the second component, independently of the third component (the cover). The parent invention is particularly useful as an intake air cleaner for internal combustion engines.\nThe present invention arose during continuing development efforts related to the noted parent invention, and provides a further mounting and retention system."}
{"text": "1. Field of the Invention\nThe present invention relates generally to recording devices used in a communications system, and, more specifically to a method for transferring video and sound data directly between a digital video recorder in a first distinct network to another digital video recorder in a second distinct network through a multichannel broadband communications network such as a cable, satellite or wireless network.\n2. Background Art\nPersonalized recording of television shows, movies, sporting events and the like has been an ever increasing trend over the last several decades. The personalized recording of these programs has migrated from the video cassette recorder (VCR) to the digital video disc (DVD) recorder to the digital video recorder (DVR) also known as a personal video recorder (PVR). While each of the DVR, DVD and VCR will provide for recording a program at one point in time for viewing at a later point in time, recordings made by DVR have distinct advantages to those recordings made by DVD and VCR.\nFirst, DVRs typically contain a computer hard drive, e.g., a disk, which allows an individual to store recorded material directly to the DVR rather than a separate medium such as a magnetic tape as is used with a VCR or compact optical disc as is used with DVDs. Second, DVRs allow an individual to pause live programming, rewind and replay a portion previously telecast. Once a particular portion of the program has been replayed and reviewed to the viewer's satisfaction, an individual may begin watching the program again at the point it was paused without missing any of the remaining program.\nA third advantage that DVRs provide is the ability to program the DVR to record at a specified time from remote locations. While individuals may program a VCR or DVD to record a particular program, generally, the VCR and DVD are programmed on location and set to a specific time to record that usually cannot be changed once the individual leaves the premises. Typically, an individual may program each of these recording means to record a specified program at a specified time. However, with DVRs, individuals who may have forgotten to program their DVR previously to record a program may direct their DVR to record a show with the use of their cellular telephone or computer from a remote location through the internet.\nAlthough the DVR has revolutionalized personal recording of television programs and movies, there is, however, one advantage that the DVD and VCR still maintain over a DVR. Recordings made by DVD and VCR may be preserved on mass storage devices such as a compact disc and magnetic tape, respectively. The disc and tape may be easily stored and moved. This allows for easy transport and viewing of the recording at almost any locale as long as there is a VCR or DVD player on location. As stated above, DVRs store recordings on a mass storage device or hard drive that is internal to the DVR, there may be no external media such as a tape or compact disc to capture the program. Although possible, it is very cumbersome to move a DVR for viewing of the program elsewhere. A VCR tape and DVD compact disc may also provide for practically unlimited storage space as long as there is room to physically store the tape or disc. Conversely, the hard disk of the DVR typically has a finite storage space that may not be used to store a large number of previously recorded programs if new programs are to be recorded.\nA further drawback of today's DVR occurs if an individual has forgotten to record a program or, after talking with another, has found a previously aired program to be interesting and worth viewing. In this instance, if the program has not been recorded by the individual to their own DVR, the individual may have missed their opportunity to view the program altogether, e.g. a sporting event. In the past, if one had recorded the program to a VCR tape or DVD compact disc, the tape or disc could easily be provided to a second individual for their own personal viewing. A recording made by DVR is not as easy to provide to a second individual for personal viewing. The internal DVR hard drive and overall size of the DVR make it difficult and cumbersome to move to a second location for viewing by another and nearly impossible to do if the second individual does not live in the same locale as the first. While the internet and electronic mail may be used to transfer previously recorded programming to another individual, the recordings are often very large and difficult to transfer and also require internet access to do so.\nTherefore, a need exists for a method to transfer video and sound data directly between a digital video recorder in a first distinct network to another digital video recorder in a second distinct network through a multichannel broadband communications network such as a cable, satellite or wireless network."}
{"text": "1. Field of the Invention\nThe present invention is directed to an improvement in a tool changing or transport mechanism for a machine tool which mechanism enables changing the tool in the spindle of the machine tool with tools which are normally carried in an adjacent storage structure.\n2. Prior Art\nIn U.S. Pat. No. 3,999,664, a tool changer with a dual axis of rotation is described, and this tool changer had many advantages over the tool changer disclosed in U.S. Pat. No. 3,760,490. The dual axis tool changing mechanism had a transfer arm, which is rotatable on a first axis and enables moving a tool from a spindle of a machine tool to a pocket in a magazine of a storage device. In order to enable inserting the tool in the pocket and also into the spindle, the transfer arm rotates the tool about a second axis which extends along the length of the arm. Thus, during transfer the tool is turned end for end.\nWhile this device has proven extremely satisfactory in many machine tools, problems may occur when the machine tool has an increased size so that the distance between the operating axis defined by the spindle and the storage axis increases. For example, if the length of the transfer arm is increased to enable transferring the tool over a greater distance, the arc of the path of movement of the tool is increased and in some instances, increases problems concerning providing sufficient clearance for movement of the tool during a tool transfer operation."}
{"text": "In recent years, increasingly higher strength has been demanded from steel plate which is used for automobiles, buildings, etc. For example, high strength cold rolled steel plate with an ultimate tensile strength of 900 MPa or more is being rapidly applied as bumpers, impact beams, and other reinforcing members. However, at the time of application of high strength steel plate, it is necessary to solve the problem of prevention of delayed fracture.\n“Delayed fracture” is the phenomenon of sudden fracture of a steel member (for example, PC steel wire, bolts) on which a high stress acts under the conditions of use. It is known that this phenomenon is closely related to the hydrogen which penetrates the steel from the environment.\nAs a factor greatly affecting delayed fracture of steel members, the steel plate strength is known. Steel plate is more resistant to plastic deformation and fracture the higher the strength, so there is a high possibility of use in an environment in which a high stress acts.\nNote that, if using a low strength steel member for a member on which a high stress acts, the member plastically deforms and fractures, so delayed fracture does not occur.\nIn a steel member which is shaped from steel plate such as steel plate for automobile use, the residual stress which occurs after shaping becomes larger the higher the steel plate strength, so there is a high concern over the occurrence of delayed fracture. That is, in a steel member, the higher the strength of the steel, the higher the concern over the occurrence of delayed fracture.\nIn the past, much effort has been made in the fields of steel bars or thick-gauge steel plate to develop steel materials taking delayed fracture resistance into consideration. For example, in steel bars and steel for bolt use, development has focused on formation of tempered martensite. It has been reported that Cr, Mo, V, and other elements which raise the temper softening resistance are effective for improvement of the delayed fracture resistance (for example, see NPLT 1).\nThis is art for causing the precipitation of alloy carbides, which act as trap sites of hydrogen, so as to change the mode of delayed fracture from grain boundary fracture to intragranular fracture.\nHowever, the steel which is described in NPLT 1 contains 0.4% or more of C and a large amount of alloy elements, so the workability and weldability which are required from steel sheet deteriorate. Further, to cause the precipitation of alloy carbides, several hours or more of heat treatment is necessary, so the art of NPLT 1 had the problem of manufacturability of steel.\nPLT 1 describes using oxides mainly comprised of Ti and Mg to prevent the occurrence of hydrogen defects. However, this art covers thick steel plate and considers delayed fracture after large heat input welding, but both the high workability and delayed fracture resistance which are demanded from steel sheet are not considered.\nIn steel sheet, since the thickness is small, even if hydrogen penetrates it, it is released in a short time. Further, in terms of workability, steel plate with an ultimate tensile strength of 900 MPa or more had almost never been used before, so the problem of delayed fracture had been treated as small. However, today, use of high strength steel sheet is rising, so development of high strength steel plate with excellent hydrogen embrittlement resistance has become necessary.\nUp to now, the art for raising the hydrogen embrittlement resistance almost all relates to steel material which is used at the proof stress or yield stress or less as bolts, steel bars, thick steel plate, and other such products. That is, the prior art is not art covering steel materials (steel plate) such as for members of automobiles where workability (cuttability, press formability, etc.) and, simultaneously, hydrogen embrittlement resistance are sought.\nUsually, a member obtained by shaping steel plate has residual stress remaining inside of the member. Residual stress is local, but sometimes exceeds the yield stress of the material steel plate. For this reason, steel plate free of hydrogen embrittlement even if high residual stress remains inside the member has been sought.\nRegarding the delayed fracture of steel sheet, for example, NPLT 2 reports about the aggravation of delayed fracture due to work-induced transformation of retained austenite. This considers the shaping of steel sheet. NPLT 2 describes an amount of retained austenite not causing deterioration of the delayed fracture resistance.\nThat is, the above report relates to high strength steel sheet which has a specific structure. This cannot be said to be a fundamental measure for improvement of the delayed fracture resistance.\nPLT 2 describes steel plate for enamelware use which is excellent in fishscale resistance as steel sheet considering hydrogen trapping ability and shapeability. This traps the hydrogen which penetrates steel plate at the time of production as oxides in the steel plate and suppresses the occurrence of “fishscale” (surface defects) which occur after enameling.\nHowever, with the art of PLT 2, the steel plate contains a large amount of oxides inside of it. If oxides disperse in the steel plate at a high density, the shapeability deteriorates, so it is difficult to apply the art of PLT 2 to steel plate for automobile use from which a high shapeability is required. Furthermore, the art of PLT 2 does not achieve both high strength and delayed fracture resistance.\nTo solve these problems, steel plate in which oxides are precipitated has been proposed (for example, see PLT 3). In such steel plate, the oxides which are dispersed in the steel plate act as trap sites which trap the hydrogen which has penetrated the steel, so dispersion or concentration of hydrogen at locations where stress concentrate and locations where delayed fracture is of a concern is suppressed.\nHowever, to obtain such an effect, steel plate must have oxides dispersed in it at a high density. Strict control of the production conditions is necessary.\nRelating to high strength steel plate, for example, there are the arts of PLTs 4 to 9. Further, relating to hot dip galvanized steel plate, for example, there is the art of PLT 10, but as explained above, it is extremely difficult to develop high strength steel plate wherein both delayed fracture resistance and good shapeability are achieved.\nPLT 11 discloses ultrahigh strength steel strip which has a tensile strength of 980 N/mm2 or more and is excellent in durability. In this ultrahigh strength steel strip, hydrogen delayed cracking resistance is considered, but basically martensite is used to handle the delayed fracture resistance (conventional method), so the shapeability is insufficient.\nPLT 12 discloses high strength steel strip which has a tensile strength of 980 MPa or more and is excellent in hydrogen embrittlement resistance. PLT 13 discloses high strength cold rolled steel plate which is excellent in workability and hydrogen embrittlement resistance.\nHowever, in all of this steel plate, the amount of particles which precipitate inside the grains is large. The hydrogen embrittlement resistance does not reach the level which is currently sought. Therefore, development of high strength steel plate which achieves both delayed fracture resistance and good shapeability has been strongly sought."}
{"text": "Landing an aircraft is one of the most challenging tasks a pilot undertakes. During the landing process, the aircraft is transitioned from operating in three dimensions of motion to operating in two dimensions of motion, and brought to a stop. To perform the landing properly, the aircraft approaches the runway within certain attitude, course, speed, and rate of descent limits. The course limits include, for example, both lateral limits and glide slope limits. An approach outside of these limits can result in the aircraft landing short of the runway, overrunning the runway, or otherwise contacting the runway surface in an uncontrolled manner. Any one of these events is undesirable.\nIn some instances visibility may be poor during approach and landing operations, resulting in what is known as instrument flight conditions. During instrument flight conditions, pilots rely on instruments, rather than visibility, to navigate the aircraft. Many airports and aircraft include instrument landing systems (ILS) to help guide aircraft during approach and landing operations. These systems allow for the display of a lateral deviation indicator to indicate aircraft lateral deviation from the approach course, and the display of a glide slope indicator to indicate vertical deviation from the glide slope.\nThe glide slope indicator is typically displayed on an aircraft primary flight display. Although the manner in which the glide slope indicator is displayed may vary, more recent glide slope indicators are configured to display a dot, or other geometric symbol, on one side of the primary flight display. The geometric symbol may be displayed relative to another geometric symbol or set of symbols, to indicate whether the aircraft is above, below, or on the desired glide slope.\nAlthough present glide slope indicators, including those described immediately above, are generally safe, reliable, and robust, these glide slope indicators do suffer certain drawbacks. For example, these glide slope indicators are displayed in a manner that may not intuitively depict the current aircraft glide slope deviation to the pilot. As a result, the pilot may not correctly maneuver the aircraft, at least initially, to reduce the glide slope deviation.\nHence, there is a need for a system and method of displaying aircraft glide slope to an aircraft flight crew that is more intuitive than presently used indicators. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description of the invention and the appended claims, taken in conjunction with the accompanying drawings and this background of the invention."}
{"text": "This invention relates generally to a combination weighing machine and, more particularly, to a product feeding and weighing structure for such a machine.\nIt is a common practice to increase the number of weigh hoppers for improving the accuracy of combination weights in a combination weighing machine. However, it is not always necessary to provide each weigh hopper with a feed hopper for feeding product to the weigh hopper since each weigh hopper is not emptied every combination cycle. Therefore, it has been proposed to associate a single feed hopper with a plurality of weigh hoppers. Japanese utility model opening gazette Nos. 60-17419 and 60-17420 disclose such product feeding and weighing structures. In the disclosed structures, each of the weighing heads includes a pair of weigh hoppers arranged side-by-side and a common feed hopper disposed above these weigh hoppers. The feed hopper is arranged to move between the pair of weigh hoppers to feed product selectively to either one which has been unloaded and is thus empty.\nHowever, the disclosed structures have disadvantages. For example the increased number of weigh hoppers also increases the number of expensive delicate weighing units. Also the distance between the weigh hoppers of each pair requires a substantial travel time for the feed hopper, which results in undesirable reduction in the machine operation speed."}
{"text": "Charge-transporting thin-films which contain a charge-transporting oligomer or polymer are commonly used as the charge injection layer, charge transport layer or emissive layer in organic electroluminescence devices (abbreviated below as “organic EL devices”).\nThese charge-transporting thin-films are required to have a high flatness. This is because surface irregularity of the charge-transporting thin-film leads to, for example, uneven light emission by the device, a shortened device life, dark spot formation, and lowering of the device characteristics due to electrical shorting, and is thought to bring about a decline in yield during the production of organic EL devices.\nAccordingly, there exists a strong desire for a technology that is industrially feasible and is capable of forming charge-transporting thin-films of high flatness.\nCharge-transporting thin-film forming processes are broadly divided into dry processes such as vacuum deposition and wet processes such as spin coating.\nComparing dry processes and wet processes, the latter processes are capable of easily producing thin-films having a high flatness and thus can be expected to improve the device yield.\nFor reasons such as this, it is wet processes that are predominantly used in the industrial fabrication of organic EL devices.\nIn light of the above, the inventors have developed charge-transporting varnishes for producing, by wet processes, charge-transporting thin-films applicable to various types of electronic devices (see, for example, Patent Documents 1 to 3).\nHowever, given the advances being made in the area of electronic displays, such as the trend in recent years toward larger devices for displaying graphics, there has been a growing need for charge-transporting varnishes which are capable of very precise and reproducible film formation over a large surface area. Existing varnishes have not always been able to fully satisfy such a need."}
{"text": "Passenger vehicles include a powertrain that is comprised of an engine, multi-speed transmission, and a differential or final drive. The multi-speed transmission increases the overall operating range of the vehicle by permitting the engine to operate through its torque range a number of times. The number of forward speed ratios that are available in the transmission determines the number of times the engine torque range is repeated. Early automatic transmissions had two speed ranges. This severely limited the overall speed range of the vehicle and therefore required a relatively large engine that could produce a wide speed and torque range. This resulted in the engine operating at a specific fuel consumption point during cruising, other than the most efficient point. Therefore, manually-shifted (countershaft transmissions) were the most popular.\nWith the advent of three- and four-speed automatic transmissions, the automatic shifting (planetary gear) transmission increased in popularity with the motoring public. These transmissions improved the operating performance and fuel economy of the vehicle. The increased number of speed ratios reduces the step size between ratios and therefore improves the shift quality of the transmission by making the ratio interchanges substantially imperceptible to the operator under normal vehicle acceleration.\nIt has been suggested that the number of forward speed ratios be increased to six or more. Six-speed transmissions are disclosed in U.S. Pat. No. 4,070,927 issued to Polak on Jan. 31, 1978; and U.S. Pat. No. 6,422,969 issued to Raghavan and Usoro on Jul. 23, 2002.\nSix-speed transmissions offer several advantages over four- and five-speed transmissions, including improved vehicle acceleration and improved fuel economy. While many trucks employ power transmissions having six or more forward speed ratios, passenger cars are still manufactured with three- and four-speed automatic transmissions and relatively few five or six-speed devices due to the size and complexity of these transmissions.\nSeven-speed transmissions are disclosed in U.S. Pat. No. 6,623,397 issued to Raghavan, Bucknor and Usoro. Eight speed transmissions are disclosed in U.S. Pat. No. 6,425,841 issued to Haka. The Haka transmission utilizes three planetary gear sets and six torque transmitting devices, including two brakes and two clutches, to provide eight forward speed ratios and a reverse speed ratio. One of the planetary gear sets is positioned and operated to establish two fixed speed input members for the remaining two planetary gear sets. Seven-, eight- and nine-speed transmissions provide further improvements in acceleration and fuel economy over six-speed transmissions. However, like the six-speed transmissions discussed above, the development of seven-, eight- and nine-speed transmissions has been precluded because of complexity, size and cost."}
{"text": "Sino-Tibetan based languages, such as Chinese, are vastly different than Latin based languages such as English. The Chinese language does not contain an alphabet. Instead, the Chinese language comprises more than 60,000 individual characters. Each of the 60,000 characters has a different meaning. Knowledge of about 1,200 characters is sufficient to read a Chinese newspaper. Chinese college graduates know about 3,000 characters.\nChinese also differs from Latin based languages in the concept of a word. In Chinese, strings of characters do not contain spaces and the interpretation of where one word ends and another starts is entirely based on context. Chinese characters are very precise in meaning, pronunciation, and in the way they are written. If a Chinese character has characters added to it in a string, the meaning of the first character is enhanced, but normally it is not changed.\nChinese characters are always pronounced as a single syllable. There are no two-syllable Chinese characters. Each Chinese character has one of five fundamental sounds. These five fundamental sounds give a singing quality to Chinese because some characters are pronounced with high tones, some with low tones, and some with tones that are rising or falling. Tone is fundamental to the language and Chinese would not be readily understood without the tones. For example, the character “ma” can either mean “mother” or “horse” or a “question” depending the tone. In China many dialects are spoken. Spoken words are almost unintelligible from one dialect to the next. However, there is only one written Chinese. Written Chinese is understood by all dialects. Other Sino-Tibetan languages such as Japanese, Korean, and Vietnamese use several characters common to Chinese. However, these languages have no common written or spoken meaning, similar to the manner in which English, Spanish, and French use a common alphabet but are not otherwise interchangeable.\nFollowing the Chinese Communist revolution in 1949, the Communist party made several changes to the Chinese language. First, the traditional method of writing Chinese from “top to bottom” and “right to left” was abandoned. The Peoples' Republic of China (PRC or mainland China) now follows Western languages and is written from “left to right” and then “top to bottom.” Second, a single dialect was chosen, Mandarin, which is now taught in all schools as the primary Chinese language. Third, the PRC altered about one quarter of the characters to reduce them to around seven lines or strokes. This form of Chinese is called “Simplified Chinese.” In the PRC, Simplified Chinese is now widely used, but the Republic of China (ROC or Taiwan) and Hong Kong still use the more elaborate form of Chinese called “Traditional Chinese.” The PRC also adopted the Hindu-Arabic numbering system used by most Western countries and the advent of the Internet is causing English to appear in many Chinese sentences.\nThe PRC also introduced “Pin Yin,” a phonetic version of Chinese to help young children learn the language. Pin Yin uses the 26 letters of the English alphabet plus 4 accents over certain vowels to indicate how the character should be pronounced. Pin Yin is normally used from about 4 years of age until around 7 years of age when the students are taught to use Chinese Characters. Pin Yin is also very helpful for tourists and businessmen to speak Chinese from phrase books. Additionally, Pin Yin is popular with computer users as it is the easiest way to enter Chinese characters from a keyboard.\nIn the computer, all Sino-Tibetan languages are represented by 16-bit characters, while English and the other Latin languages are normally represented by 8-bit characters. Traditionally, separate encodings were produced for each of the languages. English uses a 7 bit encoding called ASCII. ASCII encoding is included as the first seven bits of all the other encodings. European languages are normally 8 bit encodings and make use of the eighth bit for their special characters. Simplified Chinese uses GB2312 encoding and Traditional Chinese uses Big 5 encoding. A computer using Big 5 encoding cannot read computer code in GB2312. This multiplicity of encodings is confusing and there is no standardization between the different encodings. The Unicode consortium has developed a single encoding that incorporates all the major languages of the world. There is a strong movement to use Unicode and replace all the other encodings in computer applications. Unicode uses 16 bits for each character inside the computer. Unicode has 65,000 different characters and each of the major languages is mapped into a different section of this Unicode range. Consequently, Unicode can be used as a single encoding scheme for all of the world's languages.\nChinese characters are encoded entries which can be displayed in different font sizes. In other words, a computer may display the Chinese characters in different sizes similar to the method by which a computer displays English characters and words in different font sizes using ASCII. Changing the font size is very beneficial to students studying Chinese because the students may see the Chinese characters in greater detail.\nIndividual characters, letters, or symbols can be represented using different schemes within Unicode. Two of the most popular encoding schemes are UTF-8 and UCS-2. UTF-8 is a byte based Unicode encoding scheme which represents each character, letter, or symbol as one, two, or three bytes, each byte being eight bits. In contrast, UCS-2 is 16 bit encoding scheme which represents each character, letter, or symbol as 16 bits or four hexadecimal digits. One hexadecimal digit is equivalent to 4 bits, and 1 byte can be expressed by two hexadecimal digits. Table 1 below displays the difference between UTF-8 and UCS-2.\nTABLE 1UCS-2 (Hexadecimal)UTF-8 (Binary)Description0000 007F0xxxxxxxASCII0080 07FF110xxxxx 10xxxxxxUp to U+07FF0800 FFFF1110xxxx 10xxxxxx 10xxxxxxOther UCS-2A user may choose to encode using the UCS-2 scheme or the UTF-8 scheme depending on the user's expected needs. For example, when transmitting data from one location to another, or when storing data in a database, UTF-8 is the preferred encoding scheme due to the transmission efficiency and the storage efficiency inherent in variable byte stream length (i.e. 1-3 bytes, as shown in Table 1). However, when holding the same information in a memory, UCS-2 is the encoding scheme. Conversion functions between UCS-2 and UTF-8 are available as evidenced by United States Patent Application Publication 2003/0078921 entitled “Table-Level Unicode Handling in a Database Engine,” incorporated herein by reference.\nPrior to the development of Unicode, a computerized character translator between Simplified Chinese and Traditional Chinese within the same encoding was impossible because of the inability of GB2312 code to understand Big 5 code, and vice-versa. If the user desired a computer-implemented translation, multiple encodings had to be used which did not permit simultaneous display of both forms of data.\nSimilarly, the prior art translation programs have been unable to display Pin Yin with the proper accents. Typically, these programs would use pictures in the form of gifs or jpegs to represent the characters. The accented vowels indicate the proper tone and are essential to proper pronunciation of Pin Yin. One technique that uses only the ASCII characters is based on adding a number after the Pin Yin word to indicate the accent as illustrated in Table 2.\nTABLE 2NumberAccentDescriptionExamples1  Level Toneā ē ī ō ū2{acute over ( )}Rising Toneá é í ó ú3{hacek over ( )}Falling Tone,{hacek over (a)} {hacek over (e)} {hacek over (i)} {hacek over (o)} {hacek over (u)}then Rising Tone4{grave over ( )}Falling Toneà è ì ò ù5(None)No Change in Tonea e i o uThus, the prior art would display the word guó as guo2, the word mā as ma1, and so forth. The prior art hybrid version of Pin Yin is difficult for the beginning reader to understand because the reader must make a cognitive leap between the number and proper type and location of the accent. Therefore, a need exists for an automated method for translating between Simplified Chinese, Traditional Chinese, Pin Yin, and English. The need extends to a method for displaying the Pin Yin with the proper accent marks.\nMoreover, a need exists for assisting students with learning Chinese vocabulary. Chinese textbooks typically contain a plurality of chapters covering different subjects. Each subject presents twenty to thirty Chinese vocabulary words which are related to the subject. The student then uses the vocabulary words by themselves, then in conjunction with vocabulary words from previous chapters. Because of the encoding limitations, a computer implemented process for assisting in the development of both Simplified and Traditional Chinese vocabulary has not previously been developed. Therefore, a need exists in the art for a computer-implemented method for helping a student learn Simplified Chinese, Traditional Chinese, accented Pin Yin and English."}
{"text": "The preparation of aldehydes and alcohols by the addition of carbon monoxide and hydrogen across olefinic double bonds (hydroformylation) is known. The reaction is catalyzed by metals of Group VIII of the Periodic Table (IUPAC version) or compounds thereof which, under the reaction conditions, form carbonyls or hydridocarbonyls. Whereas cobalt and cobalt compounds were formerly used almost exclusively as catalysts, at present, the use of rhodium, even though it costs much more than cobalt, has increased. For this purpose, rhodium is used alone or in combination with complexing agents, for example organic phosphines. Whereas the reaction pressures required for the oxo synthesis with rhodium as the catalyst are from 25 to 30 MPa, pressures of from 1 to 5 Mpa suffice when rhodium complex compounds are used.\nRhodium catalysts have distinct advantages. They have higher activity and selectivity and, moreover, make it possible to operate the plant without problems in many respects, especially as to the carrying out of the synthesis and the removal of the products from the reactor. Finally, the classical oxo process based on cobalt catalysts can, in many cases, be converted to rhodium catalysts using the same apparatus, thereby minimizing capital costs.\nHowever, there are considerable difficulties in the loss-free--or at least substantially loss-free--removal and recovery of the rhodium, whether it is employed as catalyst with or without an additional complexing agent. After completion of the reaction, the rhodium is present in the hydroformylation product in the form of a solution of its carbonyl compound which may contain other ligands.\nFor workup, the pressure of the crude oxo product is normally reduced in several stages by initially reducing the synthesis pressure, which is about 1 to 30 MPa depending on the nature of the rhodium catalyst employed, to about 0.5 to 2.5 MPa. This releases the dissolved synthesis gas. Thereafter, it is possible to reduce the pressure to atmospheric. The rhodium is removed either immediately from the crude product or from the residue from the distillation of the crude product. The first route is followed when rhodium has been employed as the catalyst without additional complexing agent in the preceding hydroformylation stage; the second variant is applied when the rhodium catalyst contains other ligands in addition to carbon monoxide, for example phosphines or phosphites in complex linkages. It can also be used when, although the hydroformylation has been carried out with rhodium alone, a complexing agent has been added to the crude product after reducing the pressure to stabilize the rhodium. It is always necessary to take into account the fact that the noble metal is present in the crude product in a concentration of only a few ppm, and thus removal thereof must be carried out very carefully. Additional difficulties may arise due to the fact that, during the reduction in pressure, the rhodium, especially when it has been employed without a ligand, undergoes partial conversion into metallic form or forms polynuclear carbonyls. The result is the formation of a heterogeneous system which is composed of a liquid organic phase and a solid phase containing rhodium or rhodium compounds.\nThe recovery of rhodium from the products of the oxo synthesis including the residues of crude oxo products has been investigated many times. Such studies have led to the development of numerous processes, of which a few have also been used on an industrial scale.\nU.S. Pat. No. 4,400,547 relates to the hydroformylation of olefins with 2 to 20 carbon atoms in the presence of unmodified rhodium as the catalyst. After completion of the reaction, a complex-forming compound such as triphenylphosphine is added to the crude oxo product, and the aldehyde is removed by distillation. The distillation residue is subsequently treated with oxygen in order to eliminate the ligand from the complex compound again and to recover the rhodium in active form. Separation of the rhodium from the distillation residue is not possible in this procedure.\nThe removal of noble metals such as rhodium from high-boiling hydroformylation residues is also described in U.S. Pat. No. 3,547,964. To do this, the residues are treated with hydrogen peroxide in the presence of acids such as formic acid, nitric acid, or sulfuric acid. However, there are limits to the industrial application of the process because of the high cost of hydrogen peroxide and the difficulties of handling it.\nAccording to DE 24 48 005 C2, a rhodium-containing distillation residue is initially treated with acids and peroxides. Excess peroxides are subsequently decomposed by heating, and the aqueous solution containing the catalyst metal is reacted with hydrohalic acid or alkali metal halides and with tertiary phosphines and carbon monoxide, or compounds releasing carbon monoxide, in the presence of a water-soluble organic solvent. This procedure also requires the use of peroxides with the disadvantages described above, as well as the use of halogen-resistant materials.\nFinally, U.S. Pat. No. 4,390,473 describes a process for the recovery of rhodium and cobalt from a solution which has been employed as catalyst in a low-pressure oxo process. To remove the complex-bound metals, aqueous formic acid is added to the solution, and an oxygen-containing gas is passed through. This results in two phases, an organic phase, and an aqueous phase which contains the metals dissolved as formates. After the phases have been separated, it is possible to obtain cobalt and rhodium from the aqueous solution. In practice, however, the reducing action of formic acid has proven very bothersome. This property resulted in the rhodium being partially deposited in metallic form, and no longer amenable to recovery."}
{"text": "(1) Definition of Multi-Input Multi-Output (MIMO) Technology\nMIMO technology means a method for improving efficiency in data transmission and reception by using multiple transmitting antennas and multiple receiving antennas instead of a single transmitting antenna and a single receiving antenna. Namely, the MIMO technology is to increase capacity or improve throughput by using multiple antennas in a transmitter or receiver of a wireless communication system. Herein, multiple antennas will be referred to as MIMO.\nThe MIMO technology is an application version of a technology that a single message is completed by receiving data pieces from several antennas without depending on a single antenna path. The MIMO technology can improve a data transmission rate in a specific range or increase a system range for a specific data transmission rate. In this respect, the MIMO technology is a next generation mobile communication technology that can widely be used for a mobile station and a relay station. Also, the MIMO technology has received much attention as a next generation technology that can increase a transmission rate of mobile communication, which has reached the uppermost limit due to extension of data communication.\n(2) System Modeling in MIMO\nFIG. 1 is a schematic view illustrating a general MIMO communication system. As illustrated in FIG. 1, if the number of transmitting antennas is N7 and at the same time the number of receiving antennas increases to NR, channel transmission capacity increases theoretically in proportion to the number of antennas unlike that a plurality of antennas are used in only a transmitter or receiver. Accordingly, it is possible to improve a transmission rate and remarkably improve frequency efficiency. A transmission rate according to increase of channel transmission capacity can increase theoretically as much as a value obtained by multiplying a maximum transmission rate Ro, which corresponds to a case where one antenna is used, by an increase rate Ri of the following Equation 1.Ri=min(NT,NR)  [Equation 1]\nFor example, in a MIMO communication system that uses four transmitting antennas and four receiving antennas, a transmission rate that four times bigger than that of a single antenna system can be obtained. After theoretical capacity increase of the MIMO system has been proved in the middle of 1990, various technologies have been actively studied to substantially improve a data transmission rate. Some of the technologies have been already reflected in the standard of various wireless communications such as third generation mobile communication and next generation wireless LAN.\nUpon reviewing the recent trend of studies related to the MIMO system, active studies are ongoing in view of various aspects such as the study of information theoretical aspect related to MIMO communication capacity calculation under various channel environments and multiple access environments, the study of radio channel measurement and model of a MIMO system, and the study of time space signal processing technology for improvement of transmission reliability and transmission rate.\nIn order to describe a communication method in a MIMO system in more detail, mathematical modeling of the communication method can be expressed as follows. As illustrated in FIG. 1, it is assumed that N7 transmitting antennas and NR receiving antennas exist. First of all, a transmitting signal will be described. If there exist N7 transmitting antennas, since the number of maximum transmission information is N7, the transmission information can be expressed by a vector shown in Equation 2 as follows.s==└s1,s2, . . . ,sNT┘T  [Equation 2]\nMeanwhile, different kinds of transmission power can be applied to each of the transmission information s1, s2, . . . , sNT. At this time, supposing that each transmission power is P1, P2, . . . , PNT, transmission information of which transmission power is controlled can be expressed by a vector shown in Equation 3 as follows.ŝ=[ŝ1,ŝ2, . . . ,ŝNT]T=[P1s1,P2s2, . . . ,PNTsNT]T \nAlso, ŝ can be expressed by Equation 4 below using a diagonal matrix P.\n                    s        =                                            [                                                                                          P                      1                                                                                                                                                                                                                                                                                0                                                                                                                                                                                                                P                      2                                                                                                                                                                                                                                                                                                                                                                                                                                                                                                                    ⋱                                                                                                                                                                                          0                                                                                                                                                                                                                                                                                  P                                              N                        T                                                                                                        ]                        ⁡                          [                                                                                          s                      1                                                                                                                                  s                      2                                                                                                            ⋮                                                                                                              s                                              N                        T                                                                                                        ]                                =          Ps                                    [                  Equation          ⁢                                          ⁢          4                ]            \nMeanwhile, it is considered that a weight matrix W is applied to the information vector ŝ of which transmission power is controlled, so as to obtain N7 transmitting signals x1, x2, . . . , xNT. In this case, the weight matrix serves to properly distribute the transmission information to each antenna. Such transmitting signals x1, x2, . . . , xNT can be expressed by Equation 5 below using a vector X. In this case, Wij means a weight value between the i th transmitting antenna and the j th information. W may be referred to as a weight matrix or precoding matrix.\n                                                        x              =                            ⁢                              [                                                                                                    x                        1                                                                                                                                                x                        2                                                                                                                        ⋮                                                                                                                          x                        i                                                                                                                        ⋮                                                                                                                          x                                                  N                          T                                                                                                                    ]                                                                                        =                            ⁢                                                [                                                                                                              w                          11                                                                                                                      w                          12                                                                                            …                                                                                              w                                                      1                            ⁢                                                          N                              T                                                                                                                                                                                                                    w                          21                                                                                                                      w                          22                                                                                            …                                                                                              w                                                      2                            ⁢                                                          N                              T                                                                                                                                                                                          ⋮                                                                                                                                                                                          ⋱                                                                                                                                                                                                                                                            w                                                      i                            ⁢                                                                                                                  ⁢                            1                                                                                                                                                w                                                      i                            ⁢                                                                                                                  ⁢                            2                                                                                                                      …                                                                                              w                                                      iN                            T                                                                                                                                                              ⋮                                                                                                                                                                                          ⋱                                                                                                                                                                                                                                                            w                                                                                    N                              T                                                        ⁢                            1                                                                                                                                                w                                                                                    N                              T                                                        ⁢                            2                                                                                                                      …                                                                                              W                                                                                    N                              T                                                        ⁢                                                          N                              T                                                                                                                                                            ]                                ⁡                                  [                                                                                                                                          s                            ^                                                    1                                                                                                                                                                                          s                            ^                                                    2                                                                                                                                    ⋮                                                                                                                                                                  s                            ^                                                    j                                                                                                                                    ⋮                                                                                                                                                                  s                            ^                                                                                N                            T                                                                                                                                ]                                                                                                        =                            ⁢                              W                ⁢                                  s                  ^                                                                                                        =                            ⁢              WPs                                                          [                  Equation          ⁢                                          ⁢          5                ]            \nIf there exist NR receiving antennas, receiving signals y1, y2, . . . , yNR of the respective antennas can be expressed by a vector of Equation 6 below.y=[y1,y2, . . . ,yNR]  [Equation 6]\nMeanwhile, in case of channel modeling in the MIMO communication system, channels can be classified depending on indexes of transmitting and receiving antennas. In this case, a channel that passes from the jth transmitting antenna to the ith receiving antenna will be expressed as hij. It is noted that index of the receiving antenna is prior to index of the transmitting antenna in index of hij.\nSeveral channels can be grouped into one and then can be expressed by a vector type or a matrix type. The channels of a vector type will be described below. FIG. 2 is a diagram illustrating a channel from N7 transmitting antennas to the ith receiving antenna.\nAs illustrated in FIG. 2, the channel from N7 transmitting antennas to the ith receiving antenna can be expressed as follows.hiT=└hi1,hi2, . . . ,hiNT┘  [Equation 7]\nAlso, all channels from N7 transmitting antennas to NR receiving antennas can be expressed by Equation 8 below through the matrix of Equation 7.\n                    H        =                              [                                                                                h                    1                    T                                                                                                                    h                    2                    T                                                                                                ⋮                                                                                                  h                    i                    T                                                                                                ⋮                                                                                                  h                                          N                      R                                        T                                                                        ]                    =                      [                                                                                h                    11                                                                                        h                    12                                                                    …                                                                      h                                          1                      ⁢                                              N                        T                                                                                                                                                              h                    21                                                                                        h                    22                                                                    …                                                                      h                                          2                      ⁢                                              N                        T                                                                                                                                          ⋮                                                                                                                                                              ⋱                    ⁢                                                                                                                                                                                                                                                                          h                                          i                      ⁢                                                                                          ⁢                      1                                                                                                            h                                          i                      ⁢                                                                                          ⁢                      2                                                                                        …                                                                      h                                          iN                      T                                                                                                                    ⋮                                                                                                                                          ⋱                                                                                                                                                                                          h                                                                  N                        R                                            ⁢                      1                                                                                                            h                                                                  N                        R                                            ⁢                      2                                                                                        …                                                                      h                                                                  N                        R                                            ⁢                                              N                        T                                                                                                                  ]                                              [                  Equation          ⁢                                          ⁢          8                ]            \nSince additive white Gaussian noise (AWGN) is actually added to the channels after the above channel matrix H, AWGN n1, n2, . . . , nNR added to each of the NR receiving antennas can be expressed by a vector of Equation 9 below.n=[n1,n2, . . . ,nNR]  [Equation 9]\nThe receiving signals obtained using the above equations can be expressed by Equation 10 below.\n                                                        y              =                            ⁢                              [                                                                                                    y                        1                                                                                                                                                y                        2                                                                                                                        ⋮                                                                                                                          y                        i                                                                                                                        ⋮                                                                                                                          y                                                  N                          R                                                                                                                    ]                                                                                        =                            ⁢                                                                    [                                                                                                                        h                            11                                                                                                                                h                            12                                                                                                    …                                                                                                      h                                                          1                              ⁢                                                              N                                T                                                                                                                                                                                                                                      h                            21                                                                                                                                h                            22                                                                                                    …                                                                                                      h                                                          2                              ⁢                                                              N                                T                                                                                                                                                                                                          ⋮                                                                                                                                                                                                                                      ⋱                            ⁢                                                                                                                                                                                                                                                                                                                                                                                                  h                                                          i                              ⁢                                                                                                                          ⁢                              1                                                                                                                                                            h                                                          i                              ⁢                                                                                                                          ⁢                              2                                                                                                                                …                                                                                                      h                                                          iN                              T                                                                                                                                                                            ⋮                                                                                                                                                                                                          ⋱                                                                                                                                                                                                                                                                                  h                                                                                          N                                R                                                            ⁢                              1                                                                                                                                                            h                                                                                          N                                R                                                            ⁢                              2                                                                                                                                …                                                                                                      h                                                                                          N                                R                                                            ⁢                                                              N                                T                                                                                                                                                                          ]                                    ⁡                                      [                                                                                                                        x                            1                                                                                                                                                                            x                            2                                                                                                                                                ⋮                                                                                                                                                  x                            j                                                                                                                                                ⋮                                                                                                                                                  x                                                          N                              T                                                                                                                                            ]                                                  +                                  [                                                                                                              n                          1                                                                                                                                                              n                          2                                                                                                                                    ⋮                                                                                                                                      n                          i                                                                                                                                    ⋮                                                                                                                                      n                                                      N                            R                                                                                                                                ]                                                                                                        =                            ⁢                              Hx                +                n                                                                        [                  Equation          ⁢                                          ⁢          10                ]            \n(3) Definition of Reference Signal\nWhen a mobile communication system transmits packets, since the packets are transmitted through a radio channel, signal distortion may occur during transmission of the packets. In order to normally receive the distorted signal, a receiver should find out channel information and correct distortion of the signal as much as the channel information. In order to find out the channel information, when a signal known by both the transmitter and the receiver is received, the receiver finds out the channel information using a distortion status of the received signal. The signal known by both the transmitter and the receiver will be referred to as a pilot signal a reference signal (RS). The RS includes four types as illustrated in Table 1 depending on its purpose of use.\nTABLE 1FeaturesCQI-Common RS for channel measurement. Since UE(UseCRSEquipment) determines CQI(Channel Quality Indicator),PMI(Precoding Matrix Index), RI(Rank Indicator), etcbased on the CQI-CRS, the CQI-CRS is preferablydistributed over all bands.DM-Common RS for demodulation. But DM-CRS may be used forCRSchannel measurement. Since several UEs use the DM-CRSfor channel measurement, a precoding scheme of a specificUE cannot be applied to the DM-CRS. Accordingly, ifPDSCH(Physical Downlink Shared CHannel) is precoded,signaling is required through PDCCH(Physical DownlinkControl CHannel).NDM-Dedicated RS for non-precoded demodulation.DRSPDM-Dedicated RS for precoded demodulation. Since precodingDRSis equally applied to both PDM-DRS and PDSCH, signalingis separately required for PDCCH.\nRecently, when most of mobile communication systems transmit and receive data packets using multiple antennas, a separate reference signal for each of transmitting antennas exists, and should know a channel status between each transmitting antenna and each receiving antenna to normally receive a signal.\nMeanwhile, FIG. 3A is a diagram illustrating a structure of RS of normal cyclic prefix (CP) in a 3GPP LTE, and FIG. 3B is a diagram illustrating a structure of RS of extended cyclic prefix (CP) in a 3GPP LTE. In the current 3GPP LTE, a structure of a downlink RS is as illustrated in FIG. 3A and FIG. 3B. In one resource block, a horizontal axis means a time axis while a vertical axis means a frequency axis. One subframe includes two slots. FIG. 3A corresponds to a normal CP, wherein each slot includes seven symbols. FIG. 3B corresponds to extended CP, wherein each slot includes six symbols. The extended CP is generally used under the environment that delay is long. The number of base station transmitting antenna ports which are supported is four, and resources used to transmit RS for each of the antenna ports 0˜3 are represented by 0, 1, 2 and 3.\n(4) Time/Frequency Spacing of RS\nIn the 3GPP LTE, when considering 20 MHz system, one symbol includes 2048 samples. In this case, symbol duration is 66.67 μs. 1 μs corresponds to 10−6 seconds. Meanwhile, in case of normal CP, 160 samples are used in the first one of seven symbols as CP, and 144 samples are used in the other six symbols as CP. At this time, each of the symbols can cover channel delay spread equivalent to 5.2 μs and 4.69 μs. In case of the extended CP, since six symbols include 512 samples, they can cover channel delay spread equivalent to 16.67 μs. As seen from FIG. 3A and FIG. 3B, a transmission interval between RS for same antenna ports on the frequency axis of one symbol is 6. However, since staggering is given within one slot, a subcarrier interval between RS on the frequency axis is finally 3. In the present invention, the interval will mean a subcarrier interval on the frequency axis unless specified otherwise. In this case, it is assumed that the subcarrier interval is 15 kHz.\nMeanwhile, the transmission interval of the RS and its channel estimation capability will be considered as follows using equations.\nFor example, if inverse discrete fourier transform is applied to sequence x[i]={x0, x1, x2, . . . , xN-1}, the following Equation 11 can be obtained.\n                                          X            ⁡                          [              k              ]                                =                                    ∑                              i                =                0                                            N                -                1                                      ⁢                                          x                ⁡                                  [                  i                  ]                                            ·                              ⅇ                                  j                  ⁢                                                                          ⁢                  2                  ⁢                  π                  ⁢                                                                          ⁢                  k                  ⁢                                      ⅈ                    N                                                                                      ,                  k          =          0                ,        1        ,        …        ⁢                                  ,                  N          -          1                                    [                  Equation          ⁢                                          ⁢          11                ]            \nSince RS is sent between data at constant intervals, if 0 is filled in N number of data between the respective RSs, sequence\n      y    ⁡          [      i      ]        =      {                  x        0            ,                                    0            ,            0            ,            …            ⁢                                                  ,            0                    ︸                          z          ⁢                                          ⁢          zeros                    ,              x        1            ,                                    0            ,            0            ,            …            ⁢                                                  ,            0                    ︸                          z          ⁢                                          ⁢          zeros                    ,              x        2            ,      …      ⁢                          ,              x                  N          -          1                    ,                                    0            ,            0            ,            …            ⁢                                                  ,            0                    ︸                          z          ⁢                                          ⁢          zeros                      }  is obtained. If IDFT is applied to y[i], the following Equation 12 can be obtained.\n                                          Y            ⁡                          [              k              ]                                =                                    ∑                              m                =                0                                                              N                  ⁡                                      (                                          Z                      +                      1                                        )                                                  -                1                                      ⁢                                          y                ⁡                                  [                  m                  ]                                            ·                              ⅇ                                  j                  ⁢                                                                          ⁢                  2                  ⁢                  π                  ⁢                                                                          ⁢                  k                  ⁢                                      m                                          N                      ⁡                                              (                                                  Z                          +                          1                                                )                                                                                                                                ,                                  ⁢                  k          =          0                ,        1        ,        …        ⁢                                  ,                              N            ⁡                          (                              Z                +                1                            )                                -          1                                    [                  Equation          ⁢                                          ⁢          12                ]            \nIn the Equation 12, since y[m] and x[i] are\n  {                                                                        y                ⁡                                  [                  m                  ]                                            =                              x                ⁡                                  [                  i                  ]                                                      ,                                                              m              =                                                (                                      Z                    +                    1                                    )                                ·                i                                      ,                                                              i              =              0                        ,            1            ,            …            ⁢                                                  ,                          N              -              1                                                                                                      y                ⁡                                  [                  m                  ]                                            =              0                        ,                                                otherwise            ,                                                                                   the relationship between Y[k] X[k] can be observed using the Equations 11 and 12 in such a manner that signals are repeatedly generated as much as the number of 0 as illustrated in Equation 13.\n                              Y          ⁡                      [            k            ]                          =                  {                                                                                          X                    ⁡                                          [                      0                      ]                                                        ,                                      X                    ⁡                                          [                      1                      ]                                                        ,                  …                  ⁢                                                                          ,                                      X                    ⁡                                          [                                              N                        -                        1                                            ]                                                                      ︸                                            N                ⁢                                                                  ⁢                samples                                      ,            …            ⁢                                                  ,                                                                                X                    ⁡                                          [                      0                      ]                                                        ,                                      X                    ⁡                                          [                      1                      ]                                                        ,                  …                  ⁢                                                                          ,                                      X                    ⁡                                          [                                              N                        -                        1                                            ]                                                                      ︸                                            N                ⁢                                                                  ⁢                samples                                              }                                    [                  Equation          ⁢                                          ⁢          13                ]            \nFor example, if the interval of RS is 6, six signals are repeatedly generated in accordance with the Equation 13 as illustrated in FIG. 4. In this case, if channel delay spread does not exceed 11.11 μs, signal recovery can be performed theoretically. Namely, in case of the normal CP, since the number of CP samples is 160, 144, signal recovery can be performed even in the case that the transmission interval of RS is 6 when channel delay spread of maximum 5.2 μs and 4.69 μs occurs.\nMeanwhile, in case of the extended CP, since the number of CP samples is 512, it is supposed that channel delay spread is maximum 16.67 μs. Accordingly, a problem occurs in that interference exists between symbols as illustrated in FIG. 5.\nIn order to solve the aforementioned problem, in the LTE, staggering is applied to RS of each antenna port, so that the transmission interval of the RS becomes 3. As illustrated in FIG. 6, if the transmission interval of the RS becomes 3, channel delay spread that can be covered reaches 22.22 μs, whereby signal recovery can be performed regardless of the length of the CP.\n(5) Cyclic Delay in Time Domain\nIf sequence x[i]={x0, x1, x2, . . . , xN-1} is cyclic shifted on the time axis, the relationship with the frequency axis will be observed as illustrated in Equation 4 below.\n                                                                                                                                    X                      ′                                        ⁡                                          [                                              k                        ′                                            ]                                                        =                                    ⁢                                                            ∑                                              i                        =                        0                                                                    N                        -                        1                                                              ⁢                                                                                            x                          ⁡                                                      [                            i                            ]                                                                          ·                                                  ⅇ                                                      j                            ⁢                                                                                                                  ⁢                            2                            ⁢                            π                            ⁢                                                                                          k                                ′                                                            N                                                        ⁢                            ⅈ                                                                                              ⁢                                              ⅇ                                                  j                          ⁢                                                                                                          ⁢                          2                          ⁢                          π                          ⁢                                                      u                            M                                                    ⁢                          ⅈ                                                                                                                    ,                                                                                                        ⁢                                                      N                    =                    MZ                                    ,                                                            k                      ′                                        =                    0                                    ,                  1                  ,                  …                  ⁢                                                                          ,                                      N                    -                    1                                                                                                                          =                                ⁢                                                      ∑                                          i                      =                      0                                                              N                      -                      1                                                        ⁢                                                            x                      ⁡                                              [                        i                        ]                                                              ·                                          ⅇ                                                                        j2π                          ⁡                                                      (                                                                                                                            k                                  ′                                                                +                                uZ                                                            N                                                        )                                                                          ⁢                        ⅈ                                                                                                                                ⁢                                  ⁢                                            X              ′                        ⁡                          [                                                (                                      k                    -                    uZ                                    )                                ⁢                mod                ⁢                                                                  ⁢                N                            ]                                =                                    ∑                              i                =                0                                            N                -                1                                      ⁢                                          x                ⁡                                  [                  i                  ]                                            ·                              ⅇ                                  j2π                  ⁢                                      k                    N                                    ⁢                  ⅈ                                                                    ⁢                                  ⁢                                            X              ′                        ⁡                          [                                                (                                      k                    -                    uZ                                    )                                ⁢                mod                ⁢                                                                  ⁢                N                            ]                                =                      X            ⁡                          [              k              ]                                                          [                  Equation          ⁢                                          ⁢          14                ]            \nIn other words, if\n  ⅇ      j2π    ⁢          u      M        ⁢    ⅈ  is multiplied on the frequency axis, cyclic delay can occur on the time axis.\nIf the number of transmitting antennas that can be supported for a system that includes N number of antennas is additionally provided to support M(M>N) number of transmitting antennas, there exist an LTE UE (user equipment) that can recognize the existing N number of transmitting antennas only and an LTE-A UE that can recognize M number of transmitting antennas only. In this case, in addition to RS for supporting the existing N number of transmitting antennas, transmission of additional M-N number of RSs is required. At this time, under the environment that the LTE-A UE that recognizes M number of antennas is added to the LTE UE that recognizes the existing N number of antennas without additional signaling, efficient transmission of data and RS is required."}
{"text": "This invention relates to optical pulse sources, and in particular to a return-to-zero or carrier suppressed return-to-zero optical pulse source particularly useful in long haul optical communication systems.\nPulses in the return-to-zero (RZ) modulation format are important in long haul and ultra long haul optical fiber communication systems. RZ pulses are optical pulses where the light level reaches or xe2x80x9creturns toxe2x80x9d zero for some predetermined time interval between pulses. An RZ pulse train is a continuous stream of same shaped pulses.\nRZ pulses can be shaped to mitigate the impact of fiber chromatic dispersion, nonlinearity, and polarization mode dispersion as the optical signal travels over long fiber lengths. And, by varying the shape, duty cycle, or optical phase of an RZ pulse train, the power spectrum can be tailored to a specific system channel bandwidth and efficiency requirement. For example, RZ pulses can be formed with alternate IT phase shift, and such RZ pulses are known as carrier suppressed RZ (CSRZ) pulses.\nWhen data is encoded onto an RZ pulse stream, data xe2x80x981xe2x80x99s are represented by pulses. The absence of pulses (intervals of no light of one bit duration) represent data xe2x80x980xe2x80x99s. An important aspect of the RZ modulation format is that where two data xe2x80x981xe2x80x99s come in sequence, the light returns to zero momentarily between successive xe2x80x981xe2x80x99s.\nEarly optical RZ pulse generators simply switched on and off to generate the optical pulses. This direct modulation scheme proved ineffective because intermittent powering of the laser produced an undesirable variation of the carrier wavelength on a time scale of individual pulses (chirp). Chirp causes the pulse power spectrum to vary from the optimum spectrum of an undistorted RZ shaped pulse. This variation can adversely affect bandwidth and produces signal distortion.\nMore recent schemes for producing RZ pulses have minimized chirp to tolerable levels. One such scheme was described in U.S. Pat. No. 5,477,375, Optical Soliton Generator to Korotky, et al. The \"\"375 patent describes a CW laser followed by an amplitude modulator and a phase modulator, and finally a pulse compressor to generate an optical pulse train. The same principle was used for optical pulse generation with integrated Mach-Zehnder modulators based on GaAs and enabled integration of the pulse generator and the data modulator (Griffin, et al, Integrated 10 Gb/s Chirped Return to Zero Transmitter using GaAs/AlGaAs Modulators, OFC 2001, PD-15).\nTo date, proposed solutions for the generation of chirp-free RZ or CSRZ pulse trains are still complex and costly. Accordingly, there is a need for a low cost method and apparatus for generating RZ and CSRZ pulses.\nIn accordance with the invention, a train of RZ or CSRZ pulses is produced by passing phase modulated laser light through a delay interferometer. The parameters of the phase modulation and the delay interferometer are calculated from the desired pulse train characteristics (e.g. repetition rate, RZ or CSRZ, duty cycle). A directly modulated CW laser, or a CW laser followed by a phase modulator, produces the constant amplitude, phase modulated light. The phase modulated signal is split into two paths. One signal path is delayed with respect to the other by the calculated delay. The signals are recombined in an optical coupler to produce an RZ pulse train and/or a CSRZ pulse train."}
{"text": "In the manufacture of useful multiple-layered metallic articles such as costume jewelry, for example, it is common to first form a base metal, usually tin or brass, into a desired shape as by stamping or casting, and then by the well known and understood technique of electrode position or electroplating, applying layers of metal or metal alloy to form the finished piece. Usually a layer of copper is applied first over the base metal to give the piece ductility and provide a good adhering surface for the plating of a subsequent layer or layers of material. Next, a corrosion-resistant layer of nickel frequently referred to as a \"brightening layer\" is applied over the copper layer, and then finally, followed by a \"top coat\" or outer layer of precious metal such as gold or silver. A significant disadvantage of the well known method just described is that the ions from the nickel layer migrate to the \"top coat\" surface causing many consumers who purchase and wear such costume jewelry to suffer adverse allergic reactions, i.e. it is well known that many people have skin that is sensitive to nickel. Thus, it is not uncommon for such \"sensitive\" consumers to break out with skin irritation, itching, reddening or flaking of the skin in the region where the costume jewelry comes into contact with the wearer, e.g. the ear lobe region in the case of earrings."}
{"text": "Short message service (SMS) is a commonly used communications service nowadays. Different from a voice service, a short message service is to transport information such as text between a network and a user equipment (UE) through an SMS-service center (SMS-SC).\nA conventional short message service depends on a Mobile Station International (MSI) Public Switched Telephone Network (PSTN)/Integrated Services Digital Network (ISDN) (MSISDN).\nA short message may be sent between the calling terminal and the called terminal through a circuit switched (CS) domain or a packet switched (PS) domain or an Internet Protocol (IP) multimedia subsystem (IMS) network. With the appearance of new communications services and new terminals, although many new terminals (for example, machine type communications devices, devices with IMS subscription only) do not have an MSISDN, a short message is still required for communications.\nWhen a short message is transferred in an IMS network, a called terminal is referred to as an IMS UE; when an existing calling terminal sends a short message to the IMS UE, the calling terminal first sends a short message service-submit message to a mobile switching center (MSC) or service general packet radio service (GPRS) support node (SGSN), the short message service-submit message is required to carry the MSISDN of the called terminal. Therefore, when the called terminal does not have an MSISDN, the called terminal cannot be identified, causing that a short message cannot be normally sent, that is, a terminal cannot send a short message to an IMS UE that does not have an MSISDN."}
{"text": "The present invention relates to a method and apparatus for supporting the wall of an upwardly excavated shaft without requiring workmen to approach the wall to be supported or shielded.\nHeretofore, the process for supporting the wall of an upwardly excavated shaft generally included scaffolding in the shaft, delivering frame members up into the shaft one after another and then manually putting the frame together to complete the shielding. Such a manner of shielding naturally required the workmen to carry out their labor at the risk of falling, and further involved the danger of the destruction of the roof and the side wall or vertical wall of the shaft, as well as unsanitary conditions due to the dusty and stagnant environment.\nIn order to eliminate the above mentioned defects of the old techniques, the inventors developed, as disclosed in Japanese Patent Application No. 37542/1976, corresponding to U.S. application Ser. No. 895,103, a system for carrying out upwardly directed full face cutting accompanied by shielding of the wall, which comprises the steps of upwardly cutting a predetermined length of a shaft by a rotary bit and supporting the corresponding length of thus shaped side wall by means of at least one cylindrical frame section, whereby each of such two steps are carried out repeatedly. While such system brought about remarkable improvements in shielding or supporting the wall of an upwardly directed shaft without requiring workmen to enter into the shaft, the inventors found that the shielding of the wall including the side wall or vertical wall and the roof can be more beneficially made by providing a novel side wall support frame and a roof support frame each of which will be explained in detail hereinafter."}
{"text": "This invention relates to electrochemical sensors for determining an analyte in a fluid. There are numerous circumstances in which it is desirable to detect, measure or monitor a constituent of a fluid. Electrochemical sensors may be used for this purpose and the electrochemistry may incorporate a redox-active species whose oxidation and/or reduction is monitored as a part of the analysis.\nFor instance to measure pH, WO 2005/066618 disclosed an electrochemical sensor in which the electrochemical cell contains two organic compounds which are pH-sensitive redox systems and a ferrocene compound as an internal reference which is not sensitive to pH. To measure sulfide, WO2001/063094 and WO2004/011929 described an approach in which electrochemistry is coupled through a mediator compound to sulfide which is the intended analyte. This mediator compound is present in an electrochemical cell which is exposed to the sulfide. Both in the presence and absence of the sulfide analyte, an electrochemical oxidation and reduction of the mediator compound can take place when appropriate electrical potential is applied to the electrodes. However, one of the redox reactions of the mediator compound can also be brought about through a chemical reaction with the sulfide, and when this takes place there is a measurable change to the electrochemistry. Ferrocene carboxylate and sulfonate were suggested as mediator compounds in Electroanalysis Vol. 18, pages 1658-63 (2006) and in Electrochimica Acta Vol. 52, pages 499-50 (2006). A number of ferrocene sulfonates for possible use in this way have been described in Journal of Organometallic Chemistry Vol. 692, pages 5173-82 (2007). Experimental work in this area has, however, generally been confined to laboratory experiments at ambient temperature.\nAn issue which can arise in connection with electrochemical analytical systems is the stability of the redox active species employed, in particular stability when exposed to elevated temperatures during use. Exposure to elevated temperature may, however, be unavoidable when using an electrochemical sensor to monitor an industrial process.\nOne circumstance where there is exposure to temperature arises when carrying out analysis of fluids encountered downhole in a wellbore. Analysis of downhole fluids can be an important aspect of determining the quality and economic value of a hydrocarbon formation and can be applied to save costs and increase production at many stages of oil and gas exploration and production. Some chemical species dissolved in water (for example, Cl− and Na+) do not change their concentration when moved to the surface and information about their quantities may be obtained at the surface by analysis of downhole samples and in some cases surface samples of a flow. However, the state of chemical species, such as H+ (pH=−log [concentration of H+]), CO2, or H2S may change significantly while tripping to the surface. The change occurs mainly due to a difference in temperature and pressure between downhole and surface environment. In case of samples taken downhole, this change may also happen due to degassing of a sample (seal failure), mineral precipitation in a sampling bottle, and (especially in case of H2S)—a chemical reaction with the sampling chamber. It should be stressed that pH, H2S and CO2 are among the most critical parameters for corrosion and scale assessment. Consequently it is of considerable importance to determine their downhole values and there have been proposals for analytical sensors to be used downhole even though this is a difficult environment for an analytical system.\nRedox reactions of organic compounds solubilised in surfactant micelles have been examined, in particular for biochemical analyses carried out close to ambient temperature. One instance is Ryabov et al., J. Phys. Chem., Vol. 99, 14072 (1995) which reports voltammetry studies of ferrocene and alkyl-substituted ferrocenes in surfactants, in a biochemical context where the ferrocene redox system is coupled to glucose and glucose oxidase."}
{"text": "Microelectronic manufacturing techniques have led to the miniaturization of numerous devices. Mobile phones, personal digital assistant devices, and digital cameras are very common examples of the miniaturization trend.\nOne device that has not seen the advantage of microelectronic manufacturing techniques is the printer. Commercially available printers are large compared to many of the devices they could support. For instance, it is impractical to carry a color printer for the purpose of instantly printing photographs taken with known compact digital cameras.\nA compact print head has been described in co-pending United States patent applications and hereby incorporated by cross reference:\nUSSNTitle09/575,152, now U.S.Fluidic seal for an ink jet nozzle assemblyPat. No. 7,018,01609/575,141, now U.S.Ink jet printhead having a moving nozzle withPat. No. 6,428,133an externally arranged actuator09/575,125, now U.S.Method of manufacture of an ink jet printheadPat. No. 6,526,658having a moving nozzle with an externallyarranged actuator09/575,176, now U.S.Ink jet printhead nozzle arrayPat. No. 6,328,41709/575,147, now U.S.Nozzle guard for an ink jet printheadPat. No. 6,390,591\nAlthough digital cameras are known, they require connection to a personal computer before a hardcopy image can be obtained. Generally, it is necessary to upload images from the camera to the personal computer and then print the images on a desktop printer using a proprietary software package that makes the necessary translation between the image format taken by the camera and the format required by the printer. A camera that is able to link directly to a compact printer would be more desirable."}
{"text": "Water spray nozzles capable of providing multiple spray patterns are well known. Such nozzles are useful, for example, as nozzles that may be attached to a garden hose. By selecting one of a plurality of spray patterns, the nozzle may be used to create a jet of water that may be projected to spray water a long distance, a mist spray, a gentle soaking spray, and/or other spray patterns. Such nozzles are thus useful in replacing a number of individual nozzles that generate only a single pattern.\nGenerally, multiple pattern spray nozzle assemblies include a rotary spray head with a number of different outlets. By aligning the different outlets with an internal water conduit, water is directed to the different outlets to obtain the different patterns of spray. Such nozzle assemblies suffer from various shortcomings. For example, the rotary spray head is generally at the front end of the nozzle assembly, where water exits the nozzle assembly. Accordingly, to change the spray pattern, the water supply to the spray head must be halted or else the user is doused with the water exiting the nozzle assembly since a user's hand must typically be advanced in front of the spray path of the nozzle assembly in order to perform such function. However, even when the water supply is halted, typically by releasing a spray trigger, some water remains on the spray head. Thus, a user's hand and/or glove become wet. Furthermore, the spray trigger may be accidentally activated as the user holds the nozzle assembly and attempts to rotate the spray head, dousing the user.\nAnother shortcoming is that the spray head area of the nozzle assembly frequently comes in contact with dirt and/or other debris as the nozzle assembly is placed on the ground. Because of the water that generally remains on the spray head even when the water supply to the spray head is halted, the debris becomes caked onto to the spray head. Subsequent rotation of the spray head moves the debris into the junction where the spray head joins the body of the nozzle assembly. Eventually, the presence of debris may cause binding of the rotary spray head and/or misalignment between the internal water conduit, causing difficulty in rotation and/or leakage.\nWhat is needed is a multi-pattern spray nozzle that does not require a user to rotate a spray head located at the forward portion of a nozzle assembly by grasping the spray head itself or some other device at the front portion of the nozzle assembly. It would be beneficial if the multi-pattern spray nozzle were capable of providing multiple patterns of spray without the use of a rotating joint at the front of the nozzle assembly that is subjected to debris."}
{"text": "In a digital video broadcast system, a headend component uses Quadrature Amplitude Modulator (QAM) devices or Internet Protocol (IP) services to deliver audio/video (AV) content to a set-top box as a Moving Pictures Expert Group (MPEG) stream. The MPEG stream may include “live” advertisements or targeted advertisements. The targeted advertisements are placed in a parallel QAM or an IP service that is synchronized with the “live” advertisements. Alternatively, the targeted advertisements may be stored on a data storage device that is either integrated with the set-top box, or accessible via a network communication connection. Software on the set-top box that includes a digital video recorder (DVR) determines whether to perform a switch to the targeted advertisement based on various traits, such as demographics, purchase history, observed behavior, or the like.\nMultiple System Operator (MSO) customers are developing and deploying targeted advertising systems. These systems enable a set-top box to switch from a network service to an addressable advertising service, and back to the network service during an advertising break. The term “switch” includes the definition of “switch” in the context of the American National Standard (ANSI) Society of Cable Telecommunications Engineers (SCTE) Standard ANSI/SCTE 138 (i.e., switch describes tunes (L0 switch) and splices (L1 switch)), the definition of “tune” in the context of a QAM device, and the definition of “switch” or “transition” in the context of an IP or DVR device. In the first phase of deployment, the delivery of the network service and addressable advertising services are on separate transport streams, in which the addressable advertising content is synchronized with an advertisement on the network service. Since multiple advertisement streams are simultaneously delivered to the set-top box, the set-top box selects one of the advertisements, based on received information and locally stored selection criteria, and switches at the appropriate time.\nThe addressable advertisements contain filler, as defined in ANSI/SCTE 138. The filler follows signaled switch points which mark an opportunity to switch from the network service to the addressable advertising service, and an opportunity to switch from the addressable advertising service back to the network service. The filler is properly formatted frames that represent black pictures and silent audio, and that optimize the service acquisition procedure by allowing the set-top box the time necessary to switch to and begin decoding the new service. The duration of the filler is selected to allow for the time needed for the slowest set-top box in the network to switch to and begin decoding the new service. Since the duration of the filler is lost on the switch to the addressable advertising service and on the switch back to the network service, the duration of the filler takes away from the duration of the addressable advertising content. There is a need to minimize the duration of the filler and maximize the duration for the addressable advertising content. The presently disclosed invention satisfies this demand."}
{"text": "On-line viscometers are known for determining the yield stress and plastic viscosity of non-Newtonian or Bingham type fluids, calculated from formulae for laminar flow. However, in general, such known on-line viscometers are not capable of handling the extremes or demands which exist in the mining and mineral processing environments. For example, rotational types often suffer from reliability problems especially when applied to slurries. They also have difficulty in handling coarse solid particles and high process pressures. The static force measurement type viscometers are not suited to coarse particle slurries and will give erroneous readings when entrained air is present. Similar comments apply to vibrating element type viscometers.\nThe general technique of measuring the pressure differential between two points in a laminar fluid flowing through a tube of known diameter at a known rate is well known. This is shown, for example, in U.S. Pat. Nos. 3468158, 4680957 and 4384792. A technique using a vertical tube to measure slurry viscosity is shown in Australian Patent Application No. 31949/77. U.S. Pat. No. 3520179 discloses a rheometer or viscometer for non-Newtonian fluids using a capillary U-tube type apparatus, however this apparatus does not use differential pressure measurements, and further does not appear to be suitable for on-line measurements. U.S. Pat. No. 3465573 also discloses a U-tube type arrangement, but uses a comparison between two different paths, rather than the mere U-tube, to produce a measurement.\nU.S. Pat. No. 4680957 (Dodd) would appear to be the most relevant reference to the present invention in that it discloses a viscosity level conduit to measure viscosity. However one of the main problems with this type of viscometer is the tendency for solids to settle out of solution along the conduit thus adversely affecting the accuracy of readings from the viscometer. There is no suggestion of the use of a U-tube in this specification."}
{"text": "The present invention relates to a light set arrangement simulating lighted icicles. In this invention, the lights are arranged on downwardly dangling wires to simulate icicles where the lights are generally parallel to the wires. One such arrangement employs improved light connectors with retaining clips designed to reduce damage to the insulation surrounding the electrical wires.\nChristmas lights, as well as other decorative lights, have gained increasing popularity in recent years. Such lights are used as indoor and outdoor decorations, and are typically provided as assemblies having electrical wires or segments thereof to which light sockets and bulbs are attached. These assemblies are placed on building structures, trees, or other forms to provide lighted shapes particularly in the shape of icicles.\nIn the past, simulated icicles formed from electrical wires with light bulbs attached failed to fully reflect the linearity of an icicle. This was because the light bulbs and sockets were spaced along the electrical wires at intervals but were haphazardly attached to the wires so that the bulbs were displayed in a random arrangement with respect to the electrical wires and generally not parallel to the wires. Typically, the bulbs were displayed on opposite sides of or extending at angles from the wires. This problem was further exacerbated by the failure of the sockets holding the light bulbs to be attached to the electrical wires so that they remain fixed and generally parallel to the wires. Often the light sockets and attached electrical wires were twisted and bent around straight-edged clips causing abrasion and other damage to the insulation surrounding the electrical wires.\nThere is a need for a light set arrangement and light set connectors that provide a fixed arrangement of lights to more realistically or linearly simulate icicles. There is also a need for a light socket arrangement that reduces damage to the insulation of the electrical wires. The present invention satisfies these needs."}
{"text": "A gas turbine is conventionally comprised of a compressor, a combustor, and a turbine. The turbine is coupled to the compressor in order to drive the compressor. The combustion chamber receives fuels such as a combustion gas, and a certain amount of nitrogen, to lower the flame temperature in the combustion chamber of the combustor, which makes it possible to minimize the discharge of nitrogen oxides to atmosphere. The combustion gas may be obtained by gasification, that is, oxidation of carbon products such as coal. This partial oxidation is carried in an independent unit referred to as a gasifier. Conventionally, the gas turbine is combined with an air separation unit. The air separation unit enables at least one gas stream, mostly consisting of one of the gases of air, especially oxygen or nitrogen, to be supplied from input air. To combine the air separation unit with the gas turbine, the oxygen and nitrogen produced in the air separation unit are admitted respectively into the gasifier and the combustion chamber of the combustor."}
{"text": "Wire wrap screens comprise a perforated base pipe having a plurality of spaced axially oriented rods over which a wire generally having a trapezoidal cross section is wrapped continuously and spot welded to the rods as the winding process continues. Fairly sophisticated machines such as illustrated in FIG. 2 of US 2005/0125980 are used to wrap the single wire around the axially oriented rods to make a screen that is of a particular diameter. As is also illustrated in FIG. 5 of this reference, nested wire wrap screens of different dimensions can be produced at the same time with a single wire feed going on at the same time for constructing the screens that have different outside diameters. Even when building nested screens at the same time, each screen is assembled with a single wire feed. This is a rather slow process that could require as much as 90 minutes for the fabrication of 20 feet of screen. The fabrication process involves feeding wire on a rotating hub of a base pipe with spaced axially oriented rods and spot welding the wire to the rods periodically to maintain the wire spacing on the wraps. It is the wire spacing that determines the particle size that will be retained on the wires when the finished screen is in use.\nOther techniques called pre-pack put a filtering layer under the wire wrap as illustrated in U.S. Pat. No. 5,190,102. Other manufacturing techniques for wire wrap screens involve the use of expansion to assemble the wrapped wire and base pipe together as illustrated in U.S. Pat. No. 8,096,037.\nWhat is needed and provided by the present invention is a production method for wire wrap screens that increases the manufacturing speed of such screens. This is accomplished with a feeding technique that feeds multiple wires that begin at circumferentially spaced locations so that for every revolution of the rotating hub assembly multiple windings of wire are applied and spot welded or otherwise attached. At the opposite end of the screen the individual wires are terminated at circumferentially offset locations similar to the initial end of the manufactured screen so that the ends may be transverse to the longitudinal axis to permit end sealing with caps to be welded or simply by welding the ends to the base pipe and rods that underlie the wire wrapping. These and other aspects of the present invention will be more readily apparent to those skilled in the art from a review of the detailed description of the preferred embodiment and the associated drawings while recognizing that the full scope of the invention is to be determined from the appended claims."}
{"text": "In dynamic random access memory (DRAM) integrated circuit devices, a DRAM cell array is typically arranged in rows and columns such that a particular DRAM cell is addressed by specifying its row and column within the array. A wordline connects a row of cells to a set of bitline sense amplifiers that detect the data in the cells. In a read operation, a subset of the data in the bitline sense amplifiers is then chosen, or “column-selected” for output. DRAM cells are “dynamic” in the sense that the stored data, typically in the form of charged and discharged storage capacitors, will dissipate after a relatively short period of time. Thus, in order to retain the information, the contents of the DRAM cells must be refreshed. The charged or discharged state of the storage capacitor must be reapplied to an individual memory cell in a repetitive manner. The maximum amount of time allowable between refreshing operations is determined by the charge storage capabilities of the storage capacitors that make up the DRAM cell array. DRAM manufacturers typically specify a refresh time for which data is retained in the DRAM cells.\nA refresh operation is similar to a read operation, but no data is output. The sensing of the data in the cells by the bitline sense amplifiers is followed by a restoring operation that results in the data being rewritten to the cells. The data is, thus, “refreshed”. The refresh operation is performed by enabling a wordline according to a row address, and enabling a bitline sense amplifier. In addition, the refresh operation may be carried out by operating the bitline sense amplifier without receiving an external refresh address. In this case, a refresh address counter, which is integrated in a DRAM device chip, generates a row address subsequent to receiving an external start address.\nIt is well known that data stored in memory cells of a DRAM is retained therein by refresh operations. A self-refresh operation is performed automatically in a “standby” mode to retain the data written in the memory cells of the DRAM. In known self-refresh operations, automatic burst refresh operations can be performed at the beginning of self-refresh and at the end of self-refresh operations for shortening a refresh regulation time and for securing a stable refresh of the memory cells. U.S. Pat. No. 5,583,818 granted to You et al. on Dec. 10, 1996 discloses a semiconductor memory device having an automatic burst refresh operation in a self-refresh operation.\nA simplified description of prior art self-refresh operation is now discussed with reference to FIG. 1 that shows a self-refresh operation with automatic burst refresh function found in conventional DRAM devices. Referring to FIG. 1, in response to a command signal 111, a self-refresh mode detector 113 provides a self-refresh starting signal 115 at an entry into the self-refresh mode. A burst refresh mode controller 117 provides a burst refresh control signal 119 during a “burst self-refresh” period, in response to the self-refresh starting signal 115 and a clock signal 121 fed from a clock generator 123. A self-refresh mode controller 125 provides a self-refresh control signal 127, in response to the burst refresh control signal 119 and the self-refresh starting signal 115. The clock generator 123 provides the clock signal 121 and a self-refresh row signal 129, in response to the self-refresh control signal 127 and the burst refresh control signal 119. At an exit from the self-refresh mode, in response to the command signal 111, the self-refresh mode detector 113 ceases the self-refresh starting signal 115. Therefore, the self-refresh mode controller 125 ceases the generation of the self-refresh control signal 127. A period during the self-refresh control signal 127 being generated is a “self-refresh” period. Also, the burst refresh mode controller 117 provides the burst refresh control signal 119 during another “burst self-refresh” period. In response to the burst self-refresh control signal 119 and the self-refresh control signal 127, the clock generator 123 provides the self-refresh row signal 129 during the burst refresh period, the self-refresh period and the other burst self-refresh period. The self-refresh row signal 129 is provided to an internal row address counter which in turn provides a refresh row address signal to a row address decoder to refresh cells of the wordlines of sequentially selected rows.\nThe prior art DRAM device shown in FIG. 1 performs the self-refresh operation based on a one-cell-per-bit manner and the refresh of the DRAM cells has to be performed at a relatively high frequency. Therefore, the one-cell-per-bit self-refresh may still consume power during the self-refresh operation period. It is, thus, desirable to reduce power consumption in the self-refresh operation mode in the DRAM device."}
{"text": "There are many types of known optical recording elements. In many of the elements, the mode of operation requires a recording layer having a high absorption and the recorded marks, often referred to as pits, have low optical density or high reflectivity. The high reflectivity pits are made by ablating away the high absorption recording material, usually exposing an underlying reflective layer.\nOne of the currently popular forms of optical storage of information is the compact disk or CD. Digital information is stored in the form of high optical density marks or pits on an otherwise reflective background, as read with a focused laser diode operating in the 780-835 nm spectral region. This is the exact opposite of the above described optical recording materials. In this format, the optical information is most often in the form of read only memory or ROM. Optical information is not usually recorded in real time but rather is produced by press molding. In a typical process, the optical recording substrate is first press molded with a master containing the digital information to be reproduced. The thus formed information is then overcoated with a reflective layer and then with an optional protective layer. In those areas having the deformations or pits, the optical density is higher than in those areas not having the deformations.\nIt is desirable to produce optical recording elements which, when recorded in real time, produce a record that mimics the conventional CD on read out by generating dark marks on a reflective background.\nOne recently disclosed system of this type is the so called \"Photo CD\". In this system, conventional photographic film is first processed in a conventional manner. Then, the images from the film are scanned and digitized. The digitized information is recorded in a CD readable form on an optical recording element. Images can then be played back on a CD type player into a conventional television.\nCommercially useful materials of the type described in these references have stringent requirements. The recording layer must be able to couple with incident write laser irradiation to provide features having sufficient contrast. At the same time the unmarked portion of the disk must have sufficient reflectivity to conform to the 70% CD-ROM standard (ISO/IEC 10149). The layer must also have good stability towards light, heat and humidity for acceptable shelf life. Photo CD, for example, is a consumer product and it must be capable of withstanding extreme environments. Between the time the original images are recorded on the Photo CD and the time subsequent images are recorded, the CD might be placed in strong sunlight.\nThus, there is a continuing need for optical recording materials that have the necessary optical characteristics so that they are CD compatible, can couple with incident laser irradiation to form features with sufficient contrast, meet the minimum reflectivity standard, and yet are light stable. It is to a solution to this problem that the present invention is directed.\nU.S. Pat. No. 5,426,015 discloses optical recording elements having recording layers containing tetra dyes. These dyes comprise metallized azodianions with cationic dye counterions. These elements exhibit sensitivity that could be improved."}
{"text": "A flash memory is a non-volatile electrically erasable data storage device that evolved from electrically erasable programmable read-only memory (EEPROM). The two main types of flash memory are named after the logic gates that their storage cells resemble: NAND and NOR. NAND flash memory is commonly used in solid-state drives, which are supplanting magnetic disk drives in many applications. A NAND flash memory is commonly organized as multiple blocks, with each block organized as multiple pages. Each page comprises multiple cells. Each cell is capable of storing an electric charge. Some cells are used for storing data bits, while other cells are used for storing error-correcting code bits. A cell configured to store a single bit is known as a single-level cell (SLC). A cell configured to store two bits is known as a multi-level cell (MLC). In an MLC cell, one bit is commonly referred to as the least-significant bit (LSB), and the other as the most-significant bit (MSB). A cell configured to store three bits is known as a triple-level cell (TLC). Writing data to a flash memory is commonly referred to as “programming” the flash memory, due to the similarity to programming an EEPROM.\nThe electric charge stored in a cell can be detected in the form of a cell voltage. To read an SLC flash memory cell, the flash memory controller provides one or more reference voltages (also referred to as read voltages) to the flash memory device. Detection circuitry in the flash memory device will interpret the bit as a “0” if the cell voltage is greater than a reference voltage Vref and will interpret the bit as a “1” if the cell voltage is less than the reference voltage Vref. Thus, an SLC flash memory requires a single reference voltage Vref. In contrast, an MLC flash memory requires three such reference voltages, and a TLC flash memory requires seven such reference voltages. Thus, reading data from an MLC or TLC flash memory device requires that the controller provide multiple reference voltages having optimal values that allow the memory device to correctly detect the stored data values.\nDetermining or detecting stored data values using controller-provided reference voltages is hampered by undesirable physical non-uniformity across cells of a device that are inevitably introduced by the fabrication process, as such non-uniformity results in the reference voltages of different cells that store the same bit value being significantly different from each other. The detection is further hampered by target or optimal reference voltages changing over time due to adverse effects of changes in temperature, interference from programming neighboring cells, and numerous erase-program cycles. Errors in detecting stored data values are reflected in the performance measurement known as bit error rate (BER). The use of error-correcting codes (ECCs) can improve BER to some extent, but the effectiveness of ECCs diminishes as improved fabrication processes result in smaller cell features.\nAs illustrated in FIG. 1, an MLC flash memory has four cell voltage distributions 2, 4, 6 and 8 with four respective mean target cell voltages Vtarget0 12, Vtarget1 14, Vtarget2 16 and Vtarget3 18. Such cell voltage distributions commonly overlap each other slightly, but such overlap is not shown in FIG. 1 for purposes of clarity. During a read operation, to attempt to characterize or detect the two bits of cell data (i.e., the LSB and MSB) a flash memory device (not shown) uses three reference voltages it receives from a flash memory controller (not shown): Vref0 22, Vref1 24 and Vref2 26. More specifically, the flash memory device compares the cell voltage with Vref1 24 to attempt to detect the LSB. If the flash memory device determines that the cell voltage is less than Vref1 24, i.e., within a window 28, then the flash memory device characterizes the LSB as a “1”. If the flash memory device determines that the cell voltage is greater than Vref1 24, i.e., within a window 30, then the flash memory device characterizes the LSB as a “0”. The flash memory device also compares the cell voltage with Vref0 22 and Vref2 26 to attempt to detect the MSB. If the flash memory device determines that the cell voltage is between Vref0 22 and Vref2 26, i.e., within a window 32, then the flash memory device characterizes the MSB as a “0”. If the flash memory device determines that the cell voltage is either less than Vref0 22 or greater than Vref2 26, i.e., within a window 34, then the flash memory device characterizes the MSB as a “1”.\nThe most commonly employed ECCs are hard-decoded codes, such as BCH codes. To improve BER beyond what is commonly achievable with hard-decoded ECCs, flash memory controllers may employ soft-decoded ECCs, such as low density parity check (LDPC) ECCs. Soft decoding is more powerful in correcting errors than hard decoding, but soft input information must be provided to the ECC decoding logic. The ECC decoder soft input information is commonly provided in the form of log likelihood ratio (LLR) information. Since a flash memory device conventionally only provides hard decision outputs, i.e., it characterizes each data bit that it reads as either a “1” bit or a “0” bit in the manner described above, employing soft decoding requires that the flash memory controller generate ECC decoder soft input information.\nOne method that has been employed for generating ECC decoder soft input information (e.g., LLRs) in an MLC flash memory controller involves computing a function of the reference voltages used to read the memory device and the means and variances of the four cell voltage distributions 2, 4, 6 and 8. A faster method that has been employed for generating ECC decoder soft input information in a flash memory controller has been to map a “1” bit (hard decision) to a first fixed or predetermined value and map a “0” bit (hard decision) to a second fixed or predetermined value. For example, a “1” bit may be mapped to a decimal value of “6”, and a “0” bit may be mapped to a decimal value of “−6”. Thus, such a flash memory controller's ECC decoding logic would use decimal “6” as soft input information in response to the flash memory device reading a “1” bit and would use decimal “−6” as soft input information in response to the flash memory device reading a “0” bit. The decimal value pair of “+6” and “−6” are still hard decision values in a strict mathematical sense, but it is understood that hard decision values are a special case of soft decision values.\nDespite the benefits of employing soft-decoded ECCs, a page read sometimes fails. That is, the BER is so great that ECC decoding logic is unable to correct all erroneous bits. A common method for responding to such a page read failure is known as a “retry” or “read retry.” In a read retry, the flash memory controller may adjust the reference voltage that was used to read the page and then try to read the page again using the adjusted reference voltage.\nThe LLR value obtained from a single read may not be of sufficient quality for successful decoding by the ECC decoding logic of the controller. It is known to perform multiple reads with varying reference voltages in order to obtain more information about the read channel and higher quality LLR values. Generating high quality LLR values is critical to improving the BER of the ECC decoding logic. There are many ways to generate the LLR values. LLR values can be generated on the fly as flash memory is being read or they can be generated offline and stored in LLR tables. If the LLR values are generated on the fly as data is being read from flash memory, they are forwarded to the ECC decoder as soft input information as they are generated and used by the ECC decoder to decode the read data. If the LLR values are generated offline and stored in LLR tables, the LLR values are obtained from the table as the data is read from flash memory and provided to the ECC decoder as soft input information. The ECC decoding logic then uses the LLR values from the LLR table to decode the read data. The LLR values that are contained in the LLR table are typically predefined values that are generated once and permanently loaded into the flash memory controller."}
{"text": "This invention relates to armored cable.\nArmored cable typically has a metal sheath enclosing one or more individually insulated conductors. The sheath may be formed of a helically interlocked continuous strip of metal, or of smooth or corrugated continuous metal tubing.\nThe National Electrical Code identifies two types of armored cable: \"Type MC\" and \"Type AC\". In Type AC, the insulated conductors are individually wrapped with protective paper, and a bare ground wire extends the length of, and touches the inside wall of the sheath.\nIn Type MC cable, the insulated conductors are not individually wrapped with paper; rather they are surrounded as a group with a paper or plastic insulating wrap. The ground wire is one of the individually insulated conductors and, thus, does not touch the inside wall of the sheath. While the internal configurations of Type MC and Type AC armored cables differ, they have similar external appearances (i.e., gray-colored metal).\nArmored cable is used in wiring homes and commercial buildings. During installation, appropriate lengths of armored cable are cut from a main supply (e.g., a spool, coil, or reel) and pulled or routed through the walls and ceilings of the building.\nType AC cable must (by regulation) and Type MC cable should (for maximum safety) maintain an uninterrupted conductive path along the metal sheath from one end of the cable to the other. During installation, each end of the armored cable is typically connected to a metal terminal box using a metal connector or a plastic connector with a metal pathway that connects the armor to the box. The metal of the connector touches both the exterior of the metal sheath and the terminal box, and thus provides an electrical path from sheath to ground.\nArmored cable permits orderly wiring for various applications within a building (e.g., fire protection devices, lighting, motors). For example, wires of different sizes serving different applications may be grouped within different armored cables. Because the cables typically are installed near to each other, and because the cables look identical or similar, it is difficult to distinguish the different cables which serve different applications. The particular application associated with an armored cable may or may not be indicated by attaching an identification tag, stamping a code into the outer surface of the cable's sheath, or writing (e.g., with an indelible marker) in relatively small print somewhere on the outer surface of the sheath."}
{"text": "The present embodiments relate generally to articles of footwear, and in particular to articles with cushioning provisions and methods of making such articles.\nArticles of footwear generally include two primary elements: an upper and a sole member. The upper is often formed from a plurality of material elements (e.g., textiles, polymer sheet layers, foam layers, leather, synthetic leather) that are stitched or adhesively bonded together to form a void on the interior of the footwear for comfortably and securely receiving a foot. More particularly, the upper forms a structure that extends over the instep and toe areas of the foot, along medial and lateral sides of the foot, and around a heel area of the foot. The upper may also incorporate a lacing system to adjust the fit of the footwear, as well as to permit entry and removal of the foot from the void within the upper. In addition, the upper may include a tongue that extends under the lacing system to enhance adjustability and comfort of the footwear, and the upper may incorporate a heel counter.\nThe sole member is secured to a lower portion of the upper so as to be positioned between the foot and the ground. In athletic footwear, for example, the sole member includes a midsole and an outsole. The various sole member components may be formed from a polymer foam material that attenuates ground reaction forces (i.e., provides cushioning) during walking, running, and other ambulatory activities. The sole member may also include fluid-filled chambers, members, moderators, or other elements that further attenuate forces, enhance stability, or influence the motions of the foot, for example."}
{"text": "This invention relates in general to electrical communications and more particularly to position indication.\nWhen instrumentation computers are used during field test exercises it is customary to use Global Positioning System (GPS) receivers to provide accurate positioning information and synchronizing signals to the computers. Since a GPS receiver can provide an interface with only a single computer, each computer must be provided with its own GPS receiver. This is not practical when a large number of collocated computers are used.\nIt is therefore an object of this invention to operate several collocated instrumentation computers off a single GPS receiver.\nThis and other objects of the invention are achieved by an apparatus for providing GPS positioning information to a master computer and a plurality of slave computers from only one GPS receiver. The apparatus includes positioning information providing means adapted to be coupled between the GPS receiver and the computers for providing each computer with a replica of a positioning information signal from the GPS receiver so that each computer receives all of the positioning information it would have received if it had been connected to its own GPS receiver. In addition, the apparatus includes synchronizing means adapted to be coupled between the GPS receiver and the computers for providing each computer with a replica of a synchronizing signal from the GPS receiver so that each computer receives the synchronizing signal it would have received if it had been connected to its own GPS receiver. Further, the apparatus includes means adapted to be coupled between the GPS receiver and the master computer for passing a control signal from the master computer to the GPS receiver to set up and control the state of the GPS receiver.\nThe apparatus provides GPS positioning information to a plurality of computers through the novel use of a master-slave relationship among the computers. The master computer communicates with the GPS receiver as if the slave computers were not present and sets up and controls the state of the GPS receiver. The slave computers have no control over the GPS receiver and operate in a listen-only mode. The positioning information signal and the synchronizing signal are coupled off, converted to the proper levels, buffered, reconverted, and sent to each of the slave computers. The apparatus provides indicators showing the status of the positioning information signal and the synchronizing signal. This increases user confidence in the status of the apparatus, while enabling quick diagnosis of any problems that may arise.\nThe apparatus provides significant cost savings by reducing the number of GPS receivers and associated hardware and cabling required in the deployment of test instrumentation during developmental and operational field test exercises. It also saves on manpower because a single GPS is installed rather than multiple units. It increases reliability while reducing the consumption of power by reducing the number of active elements."}
{"text": "1. Field of the Invention\nThis invention relates to haptic interfaces and, in particular, to enabling a user to selectively engage at least part of the haptic display of a haptic interface device to produce an enhanced interaction with an environment with which the user interacts using the haptic interface device.\n2. Related Art\nA “haptic interface device” provides a haptic sensation (haptic display) to a user of the haptic interface device in response to the user's interaction with an environment with which the haptic interface device is associated. “Haptic” refers to the sense of touch: haptic interface display devices thus produce sensations associated with the sense of touch, such as texture, force (e.g., frictional force, magnetic repulsion or attraction), vibration, mass, density, viscosity, temperature, moisture, or some combination of such sensations. Haptic interface devices can be embodied in a variety of different apparatus, such as, for example, apparatus for conveying force and/or vibrotactile sensation (e.g., a stylus, a movable arm, a wheel, a dial, a roller, a slider or a vibratory surface), apparatus for conveying thermal sensation (e.g., a thermally-controlled surface or air volume), and apparatus for conveying the sensation of moisture (e.g., a moisture-controlled surface or air volume). Haptic interface devices can be used in a wide variety of applications. For example, some joysticks and mice used with computers incorporate force feedback to provide a haptic display to a user of the joystick or mouse. Some paging devices are adapted to vibrate when a paging signal is received. Some toys produce vibrations as part of the interaction with the toy. These examples give an indication of the range of applications for which a haptic interface device can be used.\nIn a conventional haptic interface device, the character of the haptic display experienced by a user is determined by a haptic model that links the state of one or more aspects of the environment to the haptic sensation provided to the user. This is illustrated in FIG. 1A. A user 101 uses an environment interaction control apparatus 103 to interact with an environment 102 via an environment interaction model 104 (either directly or via a haptic model 106, as indicated in FIG. 1A by the dashed lines between the environment interaction model 104 and the environment 102, and the environment interaction model 104 and the haptic model 106). The haptic model 106 “interprets” the user interaction with the environment 102 (based on information concerning the user interaction obtained either from the environment interaction model 104 or the environment 102, as indicated in FIG. 1A by the dashed line between the environment interaction model 104 and the haptic model 106 and the dashed arrow head entering the haptic model 106 from the environment 102) to cause a haptic display apparatus 105 to produce a corresponding haptic display. The environment interaction model 104 can also cause a non-haptic display apparatus 107 to produce a non-haptic display (e.g., a visual display and/or an audio display). However, there need not necessarily be a non-haptic display (as indicated by the dashed lines between the non-haptic display apparatus 107 and the environment interaction model 104 and user 101).\nFor example, in a previous haptic interface device used for video browsing and/or editing, a knob can be rotated to advance through the frames of a video recording, a force being applied in opposition to rotation of the knob, to simulate a detent, as transition is made from one video frame to the next. The haptic model in that haptic interface device is a relationship between applied force and position of the knob (which position corresponds to a “position” within the video recording, i.e., whether a transition between video frames is occurring or not). The same haptic display is always provided to the user for any particular user interaction, e.g., a given amount of rotation of the knob always advances the video recording a corresponding number of frames and simulates passing through the corresponding number of detents.\nThe foregoing is true of many previous haptic interface devices, i.e., the same haptic display is always provided to the user for a particular user interaction (herein, for convenience, such a haptic interface device is said to exhibit “constant” haptic behavior). Such haptic interface devices have not been as versatile as may be desirable.\nFIG. 1B is a block diagram of the functional components of a system including a reconfigurable haptic interface device. (The dashed lines and arrow heads in FIG. 1B indicate interactions between components of the system which can, but need not necessarily, occur, depending on the implementation and/or configuration of the system.) In FIG. 1B, a haptic model selection 108 results in the selection of one of two haptic models 106a or 106b for use in producing a haptic display with the haptic display apparatus 105 (i.e., the haptic model selection 108 can be used to reconfigure the haptic interface device). (More than two haptic models can be used; two are shown in FIG. 1B for simplicity.) The haptic model selection 108 may occur, for example, by specifying, before the user 101 begins using the haptic interface device, which haptic model 106a or 106b is to be used in producing the haptic display (e.g., by selecting for execution an appropriate computer program or programs used to implement the haptic model 106a or 106b). The haptic model selection 108 may or may not be made by the user 101 (as indicated by the dashed line between the user 101 and the haptic model selection 108). The haptic models 106a and 106b model different haptic behavior. Thus, unlike the haptic interface device shown in FIG. 1A, in the system shown in FIG. 1B, a particular user interaction can result in different haptic displays, depending on which of the haptic models 106a or 106b is being used to produce the haptic display (i.e., the haptic interface device does not exhibit constant haptic behavior).\nIn the system shown in FIG. 1B, a haptic interface device that is not constrained by constant haptic behavior (a reconfigurable haptic interface device) is enabled by providing multiple haptic models. There are limitations to such an approach. It may be expensive to develop, construct and/or maintain multiple haptic models. It may be difficult to establish compatibility between each of multiple haptic models and other components of the haptic interface device. It may be difficult or impossible to select a new haptic model during use of the haptic interface device by the user. Even if it is possible for the user to select a new haptic model during use of the haptic interface device, switching from one haptic model to another will result in a discontinuity in the haptic display that may be disconcerting to the user.\nIt is desirable to provide a haptic interface that is not constrained by constant haptic behavior and that eliminates or ameliorates one or more of the above-indicated limitations."}
{"text": "A technique in which a transistor is formed using a semiconductor material has attracted attention. The transistor is applied to a wide range of electronic devices such as an integrated circuit (IC) or an image display device (also simply referred to as a display device). As semiconductor materials applicable to the transistor, silicon-based semiconductor materials have been widely known, but oxide semiconductors have been attracting attention as alternative materials.\nFor example, a technique for forming a transistor using zinc oxide or an In-Ga-Zn-based oxide semiconductor as an oxide semiconductor is disclosed (see Patent Literature 1 and Patent Literature 2).\nFurthermore, in recent years, demand for integrated circuits in which semiconductor elements such as miniaturized transistors are integrated with high density has risen with increased performance and reductions in the size and weight of electronic appliances."}
{"text": "1. Field of the Invention\nThe present invention relates to a cutting insert and an indexable endmill. More specifically, the present invention relates to a cutting insert that can be replaceably attached to a multifunctional endmill capable of drilling and an indexable endmill to which the cutting insert is attached.\n2. Description of the Related Art\nA number of prior-art multifunctional endmills capable of drilling have been designed. The multifunctional endmill generally includes two cutting inserts attached at a tip end of a cylindrical tool body, and is an endmill capable of performing a cutting work in both of an XY direction (lateral direction) and a Z direction (depth direction) with respect to a work material. The cutting insert attached to such a multifunctional endmill includes a peripheral cutting edge for performing cutting in the lateral direction and a face cutting edge (bottom cutting edge) for performing cutting in the depth direction.\nAs a cutting insert attached to the multifunctional endmill, for example, the invention disclosed in Japanese Patent Laid-Open No. H10-291115(1998) is known. The multifunctional endmill described in the Japanese Patent Laid-Open No. H10-291115(1998) is of a two-cutting-edge type, which uses two cutting inserts (first cutting insert and second cutting insert). The first cutting insert and the second cutting insert are different from each other, and each of them includes the face cutting edge and the peripheral cutting edge. The cutting work is performed by the multifunctional endmill to which the two types of such cutting inserts are attached.\nIn addition, as a cutting insert attached to the multifunctional endmill, the invention disclosed in Japanese Patent Laid-Open No. H11-254207(1999) is known. The multifunctional endmill described in the Japanese Patent Laid-Open No. H11-254207(1999) is of a two-cutting-edge type, which uses two cutting inserts (first cutting insert and second cutting insert). The first cutting insert and the second cutting insert are of the same type, and two of the face cutting edge and the peripheral cutting edge are paired, and two pairs of them are provided on a top surface of the cutting insert. In other words, two corners on a diagonal line can be used on the top surface of the cutting insert. On two cutting insert seat surfaces of the endmill, the cutting inserts of the same type are disposed so that different corners are used as the cutting edges involved in cutting.\nAs described above, according to the invention of the Japanese Patent Laid-Open No. H11-254207(1999), two cutting inserts of the same type to be attached to the endmill of a two-cutting-edge type can be used so that, when the cutting edge is abraded or damaged, the two cutting inserts can be exchanged to use another pair of cutting edges. As described above, since the cutting insert of the same type can be used in turn, manufacturing costs of the cutting inserts can be reduced, and further tool management can be facilitated."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of operating a split gate type of non-volatile memory cell and a semiconductor memory device having the cells, and more particularly, to a method of operating a split gate type of non-volatile memory cell overcoming the problems of the program disturbance and endurance characteristics and a method of operating a semiconductor memory device including the cells.\n2. Discussion of Related Art\nThe split gate type of non-volatile memory cell is known in U.S. Pat. No. 5,045,488 entitled xe2x80x9cmethod of manufacturing a single transistor non-volatile semiconductor devicexe2x80x9d and U.S. Pat. No. 5,029,130 entitled xe2x80x9csingle transistor non-volatile alterable semiconductor memory devicexe2x80x9d. Methods of programming memory cells having floating gate electrodes are disclosed in U.S. Pat. No. 5,659,504 to Bude et al.\nFIG. 1 illustrates the structure of a conventional split gate type of non-volatile memory cell in U.S. Pat. Nos. 5,045,488 and 5,029,130.\nA source 12 and a drain 14 are formed on a semiconductor substrate 10. Between the source 12 and drain 14 is formed a channel 16. An insulating layer 18 is formed on the source 12, channel 16 and drain 14. A floating gate 20 is formed on a predetermined portion of the insulating layer 18 on the channel 16 and drain 12. Other insulating layer 22 is formed on the floating gate 20. Another insulating layer 24 is formed to be insulated from a control gate 26. The control gate 26 is formed on a predetermined portion of the insulating layer 24 and the insulating layer 18 on the source 12 and channel 16.\nThe operation of the split gate type of non-volatile memory cell of FIG. 1 is described with reference to FIGS. 2 to 4.\nFIG. 2 illustrates a method of erasing the split gate type of non-volatile memory cell of FIG. 1. In the drawing, the source 12 and drain 14 receive the same voltage 0V, and the control gate 26 receives the voltage Vpp higher than that applied to the source 12 and drain 14. Here, an intensive coupling from the floating gate 20 to the substrate 10 and drain 14 decreases a voltage of the floating gate 20. This kind of voltage decrease allows electrons to flow from the floating gate 20 to the control gate 26 by an F-N (Fowler-Nordheim) tunneling mechanism. Accordingly, an erasing function is ascribed to the electron of the floating gate 20 moving to the control gate 26. Through the erasing operation, the floating gate 20 is charged with (+). That is, the erasing operation is performed by the voltage difference between the floating gate 20 and the control gate 26.\nFIG. 3 illustrates a method of programming the split gate type of non-volatile memory cell of FIG. 1. A threshold voltage Nth is applied to the control gate 26. A high voltage Vpp is applied to the drain 14. xe2x80x9c0xe2x80x9d voltage is applied through the source 12 to the substrate 10, and thus the programming is performed.\nIf the high voltage Vpp is applied to the drain, the potential of the floating gate 20 is raised and the channel under the floating gate 20 is turned on. The threshold voltage Vth is applied to the control gate 26 and accordingly the channel under the control gate 26 is lightly turned on. Accordingly, electrons flow from the source 12 to drain 14. These electrons are charged in the floating gate 20 via the insulating layer 18 because of the static electricity of the floating gate 20, thus performing the programming operation. Hence, the floating gate 20 is (xe2x88x92)-charged and programmed to xe2x80x9c0xe2x80x9d.\nIn other words, the programming operation is performed in such a manner that a high voltage Vpp is applied to the drain of the memory cell to thereby bring the floating gate 20 to a predetermined voltage, and a predetermined voltage (a threshold voltage Vth of a transistor made of the control gate and the channel) is applied to the control gate 26 so that hot channel electrons generated when the current flows between the source 12 and drain 14 are injected into the floating gate.\nFIG. 4 illustrates a method of reading the conventional split gate type of non-volatile memory cell of FIG. 1. A reference voltage Vref is applied to the control gate 26; 0V to the drain 14, 2V to the source 12, and 0V to the substrate 10, thus performing the reading operation.\nIf the floating gate 20 is charged with (+), the channel 16 right under the floating gate 20 is turned on. If the voltage of the control gate 26 is raised to the reference voltage Vref for turning on the channel under the control gate 26, electrons may flow from the drain 14 to the source 12, and thereby reading the data of xe2x80x9c1xe2x80x9d.\nOn the contrary, if the floating gate 20 is charged with (xe2x88x92), the channel right under the floating gate 20 is slightly turned on or off. A voltage level of the control gate 26 and the source 12 is raised to that of the read voltage so as to turn on the channel under the control gate 26. Hence, the current cannot flow through the channel, thereby reading the data of xe2x80x9c0xe2x80x9d.\nIn other words, if the floating gate 20 is charged with (+), the current is generated through the channel 16 to thereby read the data of xe2x80x9c1xe2x80x9d, and on the contrary, if the floating gate 20 is charged with (xe2x88x92), the current does not flow through the channel 16, thereby reading the data of xe2x80x9c0xe2x80x9d.\nTherefore, the data is read by checking if the current flows through the memory cell or not by applying a predetermined voltage to the source 14 and the control gate 26. For this reason, if to perform the reading operation, the channels should be formed through the control gate and floating gate so that the current flows through the cell.\nBut, in the semiconductor memory device including the thus-structured split gate type of non-volatile memory cells, a threshold voltage Vth should be applied to the word line of a selected cell; 0V to the word line of a non-selected cell; 0V to the bit line of a selected cell; power voltage Vcc to the bit line of a non-selected cell; high voltage Vpp to the drain of a selected cell; and 0V to the drain of a non-selected cell should be respectively applied in order to generate current between the source and drain, thereby performing the programming operation. Within the programming condition, if the non-selected cell including the drain commonly connected to the drain of a selected cell is erased, the high voltage Vpp is applied to the drain of a non-selected cell, the floating gate being charged with (+), the source receiving 0V, and the substrate receiving 0V. Here, the control gate and source receive 0V, but the channel is formed by the punch through phenomenon and the current flows. The electrons conducted in the channel are injected into the floating gate, and thus programmed. Consequently, an ON-cell becomes an OFF-cell, thus causing the program interference problem.\nIn addition, the semiconductor memory device including the conventional split gate-typed memory cells repeatedly performs the programming and erasing operations. In the erasing operation, electrons of the floating gate should completely go out towards the control gate. But, they are trapped within the tunneling insulating layers. Consequently, the threshold voltage increases in accordance with the increasing number of operations, thus causing the problem of endurance characteristics.\nIt is therefore an object of the present invention to provide improved methods of operating non-volatile memory devices.\nIt is another object of the present invention to provide methods of operating non-volatile memory devices that enhance the endurance characteristics of memory cells therein.\nIt is still another object of the present invention to provide methods of operating non-volatile memory devices that increase the reliability and speed of programming and erasing operations.\nIt is still a further object of the present invention to provide methods of operating non-volatile memory devices that reduce the likelihood of programming interference between selected and non-selected memory cells.\nTheses and other objects, advantages and features of the present invention are provided by preferred methods of operating non-volatile memory cells (e.g., EEPROM cells) that include the use of negative substrate biases during programming and erasing operations. In particular, a preferred method of operating a non-volatile memory cell includes the step of erasing the memory cell by withdrawing negative charge from a floating gate therein using a positive control electrode bias and a negative substrate bias. The use of a negative substrate bias increases the potential difference between the control electrode and the floating gate and this increase results in faster and more reliable erasing. The preferred method also includes the step of programming the memory cell by accumulating negative charge on the floating gate using a positive control electrode bias, a negative substrate bias and a positive drain bias. Here, the negative substrate bias is used advantageously to reduce the likelihood that non-selected memory cells will become inadvertently programmed during operations to program selected memory cells."}
{"text": "Nearly everyone with even a modicum of cooking experience has struggled to separate two bowls that are stuck. It always seems to happen at the most inconvenient time, such as when the bowl is needed immediately in order to remove something from a hot pan.\nThere are a number of reasons why stacked bowls tend to stick. Similarly sized bowls frequently become mechanically jammed together when stacked. Bowls having different temperatures frequently stick when stacked because the bowls expand or contract at different rates as the temperatures equalize. Bowls also stick together when the weight of upper bowls in a stack force lower bowls in the stack into even lower bowls. In other cases a vacuum can be created between bowls such that the suction holds the bowls together. Bowls also get stuck when different shapes of bowls are stacked or when bowls made of different materials are stacked. In short, there are a number of reasons why bowls stick or jam together but, regardless of why they stick, it is nearly always maddening when it happens.\nPeople have been known to employ a number of means in order to separate stuck bowls, some of which are extreme. The more moderate of these means range from smartly tapping the bowls on a counter in the hope that vibration will jar them loose to using an instrument, such as a knife or screwdriver, to pry them apart. It is not unusual when frangible bowls are stuck together for one or both of the bowls to be broken in an attempt to separate them. When metal or plastic bowls are stuck, efforts to separate them frequently results in bent edges or scraped surfaces.\nA prior art method of preventing stacked bowls from sticking is to place a flat disk or similar object, such as a napkin or towel, in the lower bowl before stacking the upper bowl. Although this method works, it is not the best solution to the problem of bowls sticking together. While the problem of one bowl sticking to another may have been resolved, the new problem of preventing the bowl from sticking to the object inserted between the bowls must now be overcome.\nAccordingly, what is needed in the art is a convenient device that can be used to stack bowls so they will not stick to each other or to an object inserted between them to prevent such sticking."}
{"text": "A plurality of radio communication systems such as a cell-phone system and a radio MAN (Metropolitan Area Network) are currently used. For attaining a further speeding up and large capacity of radio communication, lively discussion is continuously performed about a next generation radio communication technology.\nFor example, in a 3GPP (3rd Generation Partnership Project) being a standardization organization, there is proposed a communication standard referred to as an LTE (Long Term Evolution) enabling communication using a frequency band of 20 MHz at a maximum. Further, as a next generation communication standard of LTE, there is proposed a communication standard referred to as an LTE-A (LTE-Advanced) enabling communication using five frequency bands (namely, a frequency band of 100 MHz) of 20 MHz at a maximum (see, for example, Non-Patent Literatures 1 and 2). In the LTE-A, the number of frequency bands to be used is proposed to be dynamically changed according to traffic (see, for example, Non-Patent Literature 3).\nFurther, in a radio communication system, from one radio communication device (e.g., a mobile station) to another radio communication device (e.g., a base station) which performs allocation control of radio resources, a random access may be performed. The random access from the mobile station to the base station is performed, for example, at the time when (1) the mobile station first accesses the base station, (2) an allocation of radio resources used for data transmission is requested to the base station, and (3) synchronization is established during reception of data from the base station, and (4) synchronization is achieved with a mobile target base station during a handover.\nThe random access includes a contention based random access and a non-contention based random access (see, for example, 10. 1. 5 section of Non-Patent Literature 4, and 5. 1 section of Non-Patent Literature 5). In the case of the random access from the mobile station to the base station, in the contention based random access, the mobile station arbitrarily selects a signal sequence from among a plurality of signal sequences and transmits it to the base station as a random access preamble. In the non-contention based random access, the base station notifies the mobile station of information in which a signal sequence is specified and the mobile station transmits a signal sequence according to the notification from the base station as the random access preamble.\nNPTL1: 3GPP (3rd Generation Partnership Project), “Requirements for further advancements for Evolved Universal Terrestrial Radio Access (E-UTRA) (LTE-Advanced)”, 3GPP TR 36.913 V8.0.1, 2009-03.\nNPTL2: 3GPP (3rd Generation Partnership Project), “Feasibility study for Further Advancements for E-UTRA (LTE-Advanced)”, 3GPP TR 36.912 V9.0.0, 2009-09.\nNPTL3: 3GPP (3rd Generation Partnership Project), “The need for additional activation procedure in carrier aggregation”, 3GPP TSG-RAN WG2 #67bis R2-095874, 2009-10.\nNPTL4: 3GPP (3rd Generation Partnership Project), “Evolved Universal Terrestrial Radio Access (E-UTRA) and Evolved Universal Terrestrial Radio Access Network (E-UTRAN); Overall description”, 3GPP TS 36.300 V9.0.0, 2009-06.\nNPTL5: 3GPP (3rd Generation Partnership Project), “Evolved Universal Terrestrial Radio Access (E-UTRA) Medium Access Control (MAC) protocol specification”, 3GPP TS 36.321 V9.1.0, 2009-12.\nIncidentally, in a radio communication system capable of performing communication by using a plurality of frequency bands, the number of frequency bands to be used according to traffic as described above is considered to be changed. However, in a method as described in the Non-Patent Literature 3, after communication is started between radio communication devices (after completing a random access procedure), a procedure is freshly performed so as to use other frequency bands except the frequency band in which communication is started. In this method, in the case where it is proved that the other frequency bands are desired to be used before starting communication (for example, in the case where a transmission data amount is proved to be large), the procedure becomes inefficient."}
{"text": "There are an increasing number of telecommunication services being proposed which, if fully implemented, could use up the allocated frequency spectrum. It is important to be able to provide such services without one service interfering with the other. If a service occupies the same frequency at the same time, there is a possibility of interference between the services which results in spectrum conflicts.\nAs will be described, so-called cognitive radios, which may employ software defined radio platforms, are capable of tailoring the transmitted output from the radio. With the advent of software defined radios, it is possible to alter the modulation type, frequency and the time of transmission to guarantee the transmission will not interfere with existing signals. Such radios may be adjusted, for instance, to inhibit transmission during times in which other signals exist. The signals from the software defined radios can also be controlled to emit non-interfering modulation formats. Programs such as the Next Generation (XG) communications funded by the Defense Advance Research Projects Agency (DARPA) propose systems where from radio scene analysis, one finds spectrum holes or White space which defines where signals may be transmitted without interference with other signals. Thus, White space refers to spaces that are not occupied by a signal.\nWhile such systems create conditions fbr transmission such that the transmitted signal does not interfere with other existing signals, spectrum utilization with such techniques is somewhat limited.\nAs will be discussed, and as part of the subject invention, it has been found that there are so-called Gray spaces where signals only partially occupy the signal space. If it were possible to be able to detect not only White spaces but also Gray spaces, then the spectrum could be more fully utilized, assuming that one could transmit non-interfering signals in the White and Gray spaces.\nMoreover, by analyzing the signal space for existing signals and providing predictors as to the future behavior of these signals, one can accurately predict future White space and Gray space. This permits robust tailoring of the transmitted signals so as not to interfere either with future signals.\nMore specifically, as telecommunications equipment evolves in capability and complexity, and Multiple-Input and Multiple-Output (MIMO) and Multi-User Detection (MUD) systems push the system throughput to its limits, it is not going to be too long before cognitive radios will reach the market place (J. Mitola, Cognitive Radio: An Integrated Agent Architecture for Software Defined Radio, Ph. D. Thesis, Royal Institute of Technology, Sweden, Spring 2000; and S. Haykin, “Cognitive Radio: Brain-Empowered Wireless Communications,” IEEE J. Select. Areas Commun., vol. 23, no. 2, pp. 201-220, February 2005).\nIn fact the IEEE 802.22 Working Group (IEEE Working Group 802.22, http://grouper.ieee.org/groups/802/22/, on Wireless Regional Area Networks (“WRANs”)), has been looking to develop a standard for a cognitive radio-based PHY/MAC/air interface for use by license-exempt devices on a non-interfering basis in spectrum that is allocated to the TV Broadcast Service on Wireless Regional Area Networks (WRAN). Ad hoc groups under the Project Authorization Request (PAR) approved by the IEEE-SA Standards Board have started developing a cognitive radio-based PHY/MAC/air interface for use by license-exempt devices on a non-interfering basis in spectrum that is allocated to the TV Broadcast Service. Moreover, cognitive radios will help the commercial as well as the military communication systems, by doing away with the need for comprehensive frequency planning. It is contemplated cognitive radios will be capable of sensing their environment, making decisions on the types of signals present, learning the patterns and choosing the best possible method of transmitting the information. They will be situation aware, and capable of making decisions to ensure error-free and smooth transfer of bits between the users. Cognitive radios will be based on software defined radio (SDR) platforms and will try to understand not only what the users want but also what the surrounding environment can provide."}
{"text": "Field of the Invention\nThe present invention relates to a method for producing a useful substance using a microorganism.\nBrief Description of the Related Art\nWhen producing amino acids such as glutamic acid and alcohols such as isopropyl alcohol by fermentation, succinic acid is produced as one of several by-products (WO2013/018734 and WO2008/114721). Consequently, the yield of the objective substance is decreased and, thus, not only the cost of the raw materials for the fermentation but also the cost of purifying the product are increased, which is uneconomical."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to a semiconductor device and a method of fabrication and more particularly to a semiconductor device with improved flicker noise characteristics.\nThis U.S. non-provisional patent application claims priority under 35 U.S. §119 of Korean Patent Application 10-2006-0045709 filed on May 22, 2006 the entire contents of which are hereby incorporated by reference.\n2. Discussion of Related Art\nDevice sizes are becoming smaller and smaller in today's semiconductor manufacturing processes. Because of these size reductions, methods of enhancing the mobility of electrons and holes are being developed. One such method induces strain in channel regions of the semiconductor device. However, strained analog MOS transistors tend to exhibit deteriorating flicker noise characteristics. Even though straining technology may have the ability to enhance the mutual conductance and cut-off frequency characteristics of analog MOS transistors, it may not be the most effective method to enhance the mobility of electrons and holes. In particular, in the case of large-scale integrated circuits (LSI) which comprise both digital and analog MOS transistors to provide completely integrated functions, it may be inappropriate to apply straining technology to both digital MOS transistors and analog MOS transistors at the same time. Accordingly, there is a need for a semiconductor device that can achieve synergies by improving both operating and noise characteristics."}
{"text": "The production of products from extruded or blown plastics generates plastic scrap. It is conventional to grind or chop the scrap and mix it with virgin or base plastics so that the scrap can be reused in the production of new products. Apparatus suitable for mixing virgin and scrap materials is disclosed in U.S. Pat. No. 4,108,334.\nOne of the problems encountered in the mixing and feeding of virgin and scrap plastics is that of maintaining a selected quantity or mass of the mixture delivered to the extruder. This problem is due in large part to variations in the densities of the materials being mixed. The density of the virgin plastic is relatively easy to control because it commonly is composed of rigid, flowable pellets of substantially uniform consistency and size. However, the scrap can, and usually does, consist of particles of greatly differing size, form, and densities. As a consequence, the bulk density of the mixture and the mass thereof discharged to an extruder may vary within wide limits.\nVariations in the quantity of material delivered during a given period of time to an extruder for extrusion is undesirable for a number of reasons, among which is the inability to maintain uniformly filled spaces between adjacent vane convolutions of a helical, extruder feed screw. For example, if a mixture having a given bulk density is fed at a constant rate to fill the space between adjacent convolutions of an extruder screw, a reduction in the bulk density of the mixture can cause such space to be underfilled. This can cause non-uniform pressure at the discharge outlet of the extruder and the production of so-called gels, i.e., unmelted plastic particles in the product. The presence of gels is undesirable inasmuch as they contribute to the generation of more scrap.\nThe problems in using a mixture of virgin and recycled scrap plastics materials have been recognized heretofore and a number of solutions advanced. One proposed solution involves slaving the rate of delivery of the virgin/scrap mixture to the speed of rotation of a variable speed extruder feed screw. As the speed of the extruder feed screw changes, the speed of the scrap delivery to the system is correspondingly changed. One difficulty with this proposal is that it does not necessarily respond to variations in the bulk density of the virgin/scrap mixture, nor does it enable a consistent feed force to be maintained on the mixture as it is delivered to the extruder. Consequently, the quantity of the mixture actually delivered to the extruder may experience wide variations.\nAnother proposal to solve the problems of feeding a mixture of virgin and scrap materials to an extruder is to premix the rigid pellets and the fluffy scrap with one another and convey the mixture to an upright, conical hopper having a vertical, motor driven feed auger which forces the mixture from the hopper into the extruder inlet. This proposal requires the use of a special, electric torque control on the feed auger drive motor, such control being operable to maintain a preset voltage on the drive motor to enable the speed of the motor to vary in response to changes in its torque. This system requires not only a special motor and drive assembly, but also an expensive torque responsive device. In addition, the time and equipment required to premix the virgin and scrap materials represents substantial expense.\nAn object of this invention is to provide a method and apparatus for feeding a mixture of different density virgin and scrap plastic materials to an extruder and wherein the feed rate of the mixture is varied in response to changes in the bulk density of the mixture to maintain substantially uniform the quantity of the mixture delivered to the extruder per unit of time, thereby overcoming the disadvantages of methods and apparatus heretofore proposed for making use of a blend of virgin and scrap materials."}
{"text": "This invention relates to electrical connectors, and, more particularly, to separable electrical connectors suited for use under high-voltage conditions. Still more particularly, this invention relates to gas actuated high-voltage bushings having a contact mounted within a bore for reciprocal movement within a bushing housing.\nHigh-voltage separable connectors innerconnect sources of energy such as transformers to distribution networks or the like. The situations typically encountered in the connection and disconnection of electrical connectors and power distributions include \"loadmake\", \"loadbreak\", and \"fault closure\". Loadmake includes the joinder of male and female contact elements, one energized and the other engaged with a normal load. An arc of moderate intensity is struck between the contact elements as they approach one another and until joinder. Loadbreak includes the separation of such mated male and female contact elements, while they supply power to a normal load. Moderate intensity arcing again occurs between the contact elements from the point of separation thereof until they are somewhat removed from one another. Fault closure includes the joinder of male and female contact elements, one energized and the other engaged with a load having a fault, e.g., a short circuit condition. A substantial arcing occurs between the contact elements as they approach one another and until joinder, giving rise to the possibility of explosion and accompanying hazard to operating personnel.\nThe prior art teaches the use of materials which emit arc-quenching gas when subjected to arcing, thus adequately dissipating the moderate intensity of arcs which occur during loadmake and loadbreak. The problem situation is fault closure, in which considerably more arc-quenching gas and mechanical assistance are required to extinguish the arc. During fault closure, the gas generated pressures may be fifty times greater than such pressures during loadmake. With respect to fault closure, the prior art has relied upon the use of the arc-quenching gas to assist in accelerating the contact elements into engagement, thus minimizing arcing time.\nA typical prior art electrical connector includes a bushing well connected to the transformer, a bushing insert which contains a female contact assembly connected to the well, and an elbow connected to a distribution line and containing a male contact to join an insert female contact in the female contact assembly. Because closure of the male and female contacts can occur under activated conditions or under fault conditions, the female contact is arranged to move within the insert to hasten the closure of the male and female contacts and thus extinguish any arc created. However, it is necessary to maintain electrical continuity during the travel of the female contact assembly. The connection between such female contact assembly and the remainder of the bushing insert must be flexible so as not to impede its movement but sufficient to carry the high currents in the circuit. Typical prior art devices include a female contact which has a piston that is moveable between a first and second position. Gas pressure which is generated by arcing during fault closure accelerates the female contact toward the male contact, thus hastening contact engagement and decreasing the time duration of the are. Mechanisms for achieving these results have not always produced sufficient current paths causing the connectors to run hot, and interfering with proper operation of the distribution network and in the extreme, leading to the destruction of the bushing inserts."}
{"text": "The present invention relates to an image display apparatus such as a head mounted display (HMD), and especially to an image display apparatus which forms images in mutually different viewing angle areas in order to display a combined image of the images joined to each other.\nObservation optical systems used in HMD are desired to be compact and lightweight. The observation optical systems are also desired to be capable of displaying an image with a wide viewing angle to increase realism in the displayed image.\nAmong observation optical systems for the HMD, a so-called ocular optical system type observation optical system is comparatively easy to be miniaturized and provided with a wide viewing angle. However, such an ocular optical system type observation optical system requires, for providing a wide viewing angle with one image-forming element, a large-size image-forming element having a large image-forming area where an original image is formed.\nIn contrast thereto, a so-called tiling technology is proposed which enables displaying of an image with a wide viewing angle by using plural small-size image-forming elements having a small image-forming area. The tiling technology joins plural enlarged images to each other which are formed by light fluxes from plural original images formed on the plural image-forming elements to enable the displaying of one enlarged combined image. Such observation optical systems (HMD) for displaying an image with a wide viewing angle by the tiling technology are disclosed in Japanese Patent Laid-Open No. 09-166759 and Japanese Patent Laid-Open No. 11-326820.\nJapanese Patent Laid-Open No. 09-166759 discloses an observation optical system in which two optical units are provided for one eye, each optical unit including an image-forming element and a three-surface prism having an entrance surface, an internal total reflection/exit surface and a reflection surface, and the two optical units are disposed at upper and lower positions so as not to be overlapped with each other in a right and left direction. In this observation optical system, two enlarged images are formed by the two optical units at right and left positions, and thereby one enlarged combined image with a horizontally wide viewing angle is displayed. When a cross section where a light flux (optical path) is folded by at least twice reflections by an optical element such as the three-surface prism is defined as a local meridional cross section, the “right and left (horizontal) direction” in Japanese Patent Laid-Open No. 09-166759 corresponds to a direction along a local sagittal cross section orthogonal to the local meridional cross section.\nJapanese Patent Laid-Open No. 09-166759 also discloses an example in which the two optical units disposed at the upper and lower positions are relatively moved so as to be overlapped with each other in a visual axis direction of an observer. A region where the two optical units are overlapped with each other allows observation of an outside optical image (so-called optical see-through observation). In this example, enlarged images formed by the two optical units are overlapped with each other at a central region in the right and left direction to form one enlarged combined image.\nJapanese Patent Laid-Open No. 11-326820 discloses an observation optical system including two division optical systems which form two enlarged images, the two division optical systems using one prism having a symmetric shape in a local meridional cross section. In this observation optical system, a light flux from an original image formed on one image-forming element enters the prism through an entrance surface thereof, is totally reflected by an internal total reflection/exit surface thereof, is then reflected by a concave mirror surface, and finally exits from the prism through the internal total reflection/exit surface to proceed to an exit pupil. A light flux from an original image formed on the other image-forming element proceeds to the same exit pupil in a same manner as that of the light flux from the one image-forming element. The two concave mirror surfaces that cause the light fluxes from the two image-forming elements to proceed to the same exit pupil are closely adjacent (joined) to each other.\nHowever, the latter observation optical system disclosed in Japanese Patent Laid-Open No. 09-166759 does not contribute to miniaturization of the observation optical system (in other words, of the HMD) in a direction of the local meridional cross section (vertical direction) in which the observation optical system is decentered, because its viewing angle is not divided in the direction of the local meridional cross section while the viewing angle is divided in a direction of the local sagittal cross section (horizontal direction) to obtain a wide viewing angle. Each of the observation optical systems disclosed in Japanese Patent Laid-Open No. 09-166759 aims to enable the see-through observation through the overlapped region of the two optical units. Therefore, in the local meridional cross section, the viewing angles of the two optical units are the same, and parts thereof are overlapped with each other. Therefore, image division cannot be performed in the local meridional cross section, which is not effective for reducing the thickness or vertical size of the observation optical system.\nMoreover, the former observation optical system disclosed in Japanese Patent Laid-Open No. 09-166759 and the observation optical system disclosed in Japanese Patent Laid-Open No. 11-326820 form a joined portion of the two images formed by the light fluxes from the two original images (corresponding to a joined portion of the two concave mirror surfaces) in the enlarged combined image. The joined portion of the two concave mirror surfaces generates light scattering and flare, which makes the joined portion of the images noticeable."}
{"text": "1. Field of the Invention\nThe present invention relates to a novel method of producing a catalyst useful in the reaction of ethylene, oxygen and acetic acid in the vapor phase to form vinyl acetate. In particular, the present invention is directed to a novel method of forming a catalyst useful in the catalytic formation of vinyl acetate in which said catalyst comprises metallic palladium and gold deposited on a suitable porous support.\n2. Description of the Prior Art\nIt is known in the an that vinyl acetate can be produced by reacting ethylene, oxygen and acetic acid in the gaseous phase and in the presence of a catalyst comprising palladium, gold and an alkali metal acetate supported on certain carrier materials such as silica. Generally, such catalyst systems exhibit a high activity. Unfortunately, results utilizing such palladium and gold catalysts have been inconsistent. This inconsistency appears to be based somewhat on the distribution pattern or profile of the catalyst components which are deposited on and in relation to the support. For example, when use is made of the known vinyl acetate catalyst systems comprising a porous support with palladium and gold, the metal components deposited at or about the support interiors or central regions do not contribute significantly to the reaction mechanism, since the reactants are not readily able to diffuse into the central or inner regions of the porous network of the catalyst. Hence, the reaction occurs substantially only at the outermost or surface regions of the catalyst. The catalyst components in the interior regions of the support do not largely contribute to the reaction scheme, resulting in a reduction in catalytic efficiency per unit weight of the catalyst components. Furthermore, the use of a highly active catalyst at times gives rise to side reactions and, therefore, leads to a reduced selectivity to vinyl acetate.\nVarious patents have been granted based on the desire to more evenly distribute and anchor the gold and palladium catalytic components within a narrow band on the support surface to provide a vinyl acetate catalyst having high yield, good selectivity and long life. Examples of such patents include U.S. Pat. Nos. 4,087,622; 4,048,096; 3,822,308; 3,775,342 and British Patent 1,521,652.\nThe basic method of forming the vinyl acetate catalyst containing palladium and gold deposited on a catalyst support comprises ( 1) impregnating the support with aqueous solutions of water-soluble palladium and gold compounds, (2) precipitating water-insoluble palladium and gold compounds on the catalyst support by contacting the impregnated catalyst support with a solution of compounds capable of reacting with the water-soluble palladium and gold compounds to form the insoluble precious metal compounds (3) washing the treated catalyst with water to remove anions which are freed from the initially impregnated palladium and gold compounds during precipitation and (4) convening the water-insoluble palladium and gold compounds to the free metal by treatment with a reducing agent. A final treatment usually involves (5) impregnating the reduced catalyst with an aqueous alkali metal acetate solution and (6) drying the final catalyst product.\nPrior an attempts to provide a uniform distribution of the palladium and gold metals on the support have involved some manipulation of the above mentioned steps and/or by using support materials having various specified pore dimensions."}
{"text": "This invention is related generally to the field of aerosolization of liquids, and in particular to the aerosolization of liquids using an aerosol generator that operates at ultrasonic vibrational frequencies. More specifically, the invention relates to techniques for vibrationally isolating an aerosol generator that is connected to another structure, such as the housing of an aerosolization device, when aerosolizing a liquid.\nAerosolization of liquids is an important aspect of many businesses. For example, liquids are commonly aerosolized in connection with drug delivery, air humidification, deodorant or insecticide delivery, and the like. One exemplary way to aerosolize liquids is by supplying liquid to a vibratable member having a plurality of apertures and vibrating the vibratable member at ultrasonic frequencies. One way to vibrate the vibratable member is by use of a piezoelectric transducer. Such techniques are described in, for example, U.S. Pat. Nos. 5,164,740; 5,938,117; 5,586,550; 5,758,637 and 6,085,740, the complete disclosures of which are herein incorporated by reference.\nWhen vibrating such vibratable members at ultrasonic frequencies, there is a need to ensure that a maximum amount of the vibrational energy is transferred from the piezoelectric transducer to the vibratable member, rather than to some surrounding structure. Otherwise, performance losses or performance variations may occur because of the forces that are transmitted through the material which couples the aerosol generator to surrounding structures, such as the housing of an aerosolization device.\nHence, this invention is related to ways to maximize the amount of vibrational energy transferred to the vibratable member, thereby maximizing the efficiency of the aerosol generator. In this way, the repeatability and performance of the aerosol generator are enhanced, irrespective of the devices into which the aerosol generators are integrated.\nThe invention provides for the vibrational isolation of an aerosol generator from surrounding structures. In one embodiment, this is accomplished by the design of an aerosol generator that comprises a vibratable member having a front, a rear, an outer periphery and a plurality of apertures extending between the front and the rear. A support element is disposed about the outer periphery of the vibratable member. A vibratable element is coupled to the support element and is configured to vibrate the vibratable member at ultrasonic frequencies. An isolating structure is coupled to the support element and is configured to couple the aerosol generator to a support structure, such as the housing of an aerosolization device. The isolating structure has a vibrational mechanical impedance that is sufficient to substantially vibrationally isolate the aerosol generator from the support structure. In this way, the aerosol generator may be operated at increased efficiencies and in a repeatable manner when coupled to surrounding structures.\nConveniently, the isolating structure and the support element may be integrated into a single component, thereby facilitating its manufacture. In one aspect, the isolating structure may comprise a plurality of arms that extend from the support element. These arms may have a wide variety of shapes and contours. For example, the arms may be bent, crimped, curved, or the like to facilitate vibrational isolation.\nIn another aspect, the isolating structure may comprise one or more elastomeric or plastic members. For example, the isolating structure may be constructed of an elastomeric or plastic washer. Conveniently, the washer may be coupled to the support element by forming tabs in the support element and inserting the washer between the tabs. As another example, the isolating structure may be constructed of a plurality of discrete elastomeric members or bellows that extend from the support element.\nTo facilitate vibrational isolation, the isolating structure may be configured so that it has a resonant frequency that is outside the operating range of the aerosol generator. Such an operating range for the aerosol generator may be about 50 kHz to about 250 kHz.\nIn a further aspect, the vibratable member may be dome shaped and include tapered apertures. Examples of such vibratable members are described in U.S. Pat. Nos. 5,586,550, 5,758,637 and 6,085,740, previously incorporated by reference.\nThe invention further provides an exemplary method for aerosolizing liquids. Such a method utilizes an aerosol generator having a vibratable member with apertures and a vibratable element to vibrate the vibratable member. According to the method, liquid is supplied to the vibratable member and the vibratable element is used to vibrate the vibratable member at an ultrasonic frequency to eject liquid droplets through the apertures. During vibration, an isolating structure is used to substantially vibrationally isolate the aerosol generator to enhance the operating performance of the aerosol generator. Further, the vibratable member may be vibrated at a frequency that is different from a fundamental frequency of the isolating structure to enhance the efficiency of the aerosol generator. As previously mentioned, a variety of isolating structures may be used to vibrationally isolate the aerosol generator from any surrounding structures. Such isolating structures also have resonant frequencies outside of the operating range of the aerosol generator.\nIn a further embodiment, the invention provides a method for forming an aerosol generator. According to the method, an isolating structure is stamped or coined out of a sheet of material. A vibratable member having a plurality of apertures is coupled to the isolating structure, and a vibratable element, such as a piezoelectric transducer, is coupled to the isolating structure or the vibratable member. The vibratable element is used to vibrate the vibratable member at ultrasonic frequencies while the isolating structure is used to vibrationally isolate the aerosol generator from surrounding structures. By forming the isolating structure in this way, the cost of producing the aerosol generator may be greatly reduced and the aerosol generator may be produced in higher volumes.\nIn one aspect, the isolating structure comprises an annular member and a plurality of arms extending from the annular member. In another aspect, the arms are bent or crimped after the isolating structure has been stamped. In a further aspect, the vibratable member is coupled across a central opening of the annular member, and the vibratable element comprises an annular piezoelectric element that is coupled to the annular member.\nAnother method for forming an aerosol generator uses a support element having an outer periphery. A plurality of tabs are formed in the outer periphery of the support element. This may be accomplished by making a pair of cuts in the support element and then bending the material between the cuts away from the support element. A vibratable member having a plurality of apertures is coupled to the support element, and a vibratable element is coupled to the support element or the vibratable member and is vibratable at ultrasonic frequencies. A gasket is coupled about the support element, with the gasket being received into the tabs. The gasket has a mechanical vibrational impedance that is sufficient to substantially vibrationally isolate the aerosol generator. Such a process is useful in producing an isolated aerosol generator in a time and cost efficient manner."}
{"text": "Various implantable medical devices are advantageously inserted within various body vessels, for example to improve blood flow through a restricted or weakened vessel. Minimally invasive techniques and instruments for placement of intraluminal medical devices have been developed to treat and repair undesirable conditions within body vessels. Various percutaneous methods of implanting medical devices within the body using intraluminal transcatheter delivery systems can be used to treat a variety of conditions. One or more intraluminal medical devices can be introduced to a point of treatment within a body vessel using a delivery catheter device passed through the vasculature communicating between a remote introductory location and the implantation site, and released from the delivery catheter device at the point of treatment within the body vessel. Intraluminal medical devices can be deployed in a body vessel at a point of treatment and the delivery device subsequently withdrawn from the vessel, while the medical device retained within the vessel to provide sustained improvement in blood flow or to increase vessel patency.\nImplantable medical devices are effective for minimally invasive treatment of vascular occlusions such as atherosclerosis and restenosis and treatment of weakened or diseased vessels. Such implantable medical devices reestablish a flow lumen, reinforces the weakened vessel, and prevents occlusion or stenosis.\nCoatings are often applied to the implantable medical device to improve the biocompatibility of the device and minimize or prevent occlusion of the device. Methods of coating the abluminal surface of a tubular medical device are known. U.S. published patent application No. 2005/0233061 A1 describes a method and apparatus for coating a medical device using a coating head. In one embodiment a slide coating is applied to the abluminal surface of a medical device. U.S. published patent application No. 2005/0196518 A1 describes a method of coating a medical device by substantially simultaneously applying a coating composition and partially drying the coating composition. U.S. published patent publication No. 2005/0147734 A1 describes a method for application of therapeutic and protect coatings to the abluminal tubular medical devices by placing the medical device on a core and passing it through an extrusion coating machine.\nInhibiting or preventing thrombosis and platelet deposition on an implantable device within the body is important in promoting continued function of the medical device within the body, particularly within blood vessels. Post-implantation thrombosis and platelet deposition on surfaces of implantable medical devices prosthesis undesirably reduce the patency rate of many implantable medical devices. For example, thrombosis and platelet deposition within an endovascular prosthesis may occlude the conduit defined by the endovascular prosthesis or compromise the function of an implanted valve by limiting the motion or responsiveness of moveable portions of the device such as valve leaflets. Many factors contribute to thrombosis and platelet deposition on the surfaces of implanted prosthesis. The properties of the material or materials forming the endovascular prosthesis are believed to be one important factor that can contribute to the likelihood of undesirable levels of post-implantation thrombus formation or platelet deposition on the implanted device. The formation of blood clots, or thrombus, on the surface of an endovascular prosthesis can both degrade the intended performance of the prosthesis and even undesirably restrict or occlude desirable fluid flow within a body vessel. Coatings may be used to prevent occlusion of the implantable medical device. A non-thrombogenic coating may be used to minimize thrombosis on the blood contact surface of the device.\nNon-thrombogenic coatings are preferably applied to the luminal surface of the medical device. U.S. Pat. No. 7,112,298 describes a method of forming a medical device comprising applying a polymer coating to a mandrel and constructing the remainder of the medical device around the polymer coated mandrel. What is needed are improved methods for coating medical devices on the luminal and/or abluminal surface. Methods for coating luminal surfaces are useful forming a non-thromobogenic blood contact surface on the medical device."}
{"text": "Alkaline disinfectant and/or sanitizer solutions, are known. Effective alkaline disinfectant and/or sanitizer solutions may be especially useful in food handling and processing such as in grocery, dairy, restaurant and food service applications, where surfaces (e.g., equipment) may be vulnerable to problems with biofilm formation due to the prevalence of microbial food sources. The 2-part alkaline disinfectant and/or sanitizer solutions are especially useful on food contact surfaces in food and dairy processing.\nIn some instances, it is useful to activate a disinfectant/sanitizer cleaner solution via addition of an activator solution to produce an activated (i.e. “in-use”) disinfectant and/or sanitizer solution at or near a substrate to be treated. By combining a disinfectant/sanitizer cleaner solution with an activator solution at or near the substrate to be treated, efficacy of the resulting activated disinfectant and/or sanitizer solution may be increased, for example by keeping incompatible compounds separate until time of use. An example of one such “2-part system,” is a system comprising Ultra Disinfectant Cleaner Solution 1 and Ultra Activator Solution, each from Sterilex® (Hunt Valley, Md.). In use, Sterilex® Ultra Disinfectant Cleaner Solution 1 and Sterilex® Ultra Activator Solution may be mixed together in equal volumes and diluted with water to create an alkaline in-use disinfectant and/or sanitizer solution. The alkaline in-use disinfectant and/or sanitizer solution may be applied to a substrate to be treated using any suitable method including, but not limited to, being applied as a foam or spray, used as soaking solution and/or circulated through or around a substrate.\nOne industrial area in which 2-part systems may be particularly useful, is to treat heat exchanger coils which can be found in spiral freezers and other freezing/chilling equipment used for processing frozen foods. Additional industrial areas in which 2-part systems may be particularly useful include overhead air handling units, cooling tunnels, aluminum paneling and walls. However, these and other substrates may comprise components that may be vulnerable to corrosive effects of compounds commonly found in alkaline in-use disinfectant and/or sanitizer solutions. For example, aluminum, which is a soft metal, is commonly used as a component of coils in spiral freezers and other equipment. However, aluminum is prone to rapid corrosion upon contact with alkaline substances.\nWhile various 2-part alkaline disinfectant and/or sanitizer solutions have been made and used, it is believed that no one prior to the inventor(s) has made or used the invention described in the appended claims, which reduce or eliminate the corrosive effects of alkaline disinfectant and/or sanitizer solutions."}
{"text": "1. Field of the Invention\nThe present invention relates to an ink jet recording apparatus, an ink cartridge and an ink jet head unit, and more particularly to an ink jet recording apparatus, an ink cartridge and an ink jet head unit, which are capable of high quality color recording with good water resistance and no spreading of ink.\nThe term xe2x80x9crecordingxe2x80x9d includes application of ink to (or printing of) all kinds of ink supporting body such as cloth, threads, paper and sheet materials. This invention can be applied to any kind of equipment that use such recording materials as paper, cloth, nonwoven cloth and OHP sheets for over-head projector. Among equipment to which the invention is applicable are office equipment such as printers, copying machines and facsimiles, and also mass production equipment.\n2. Description of the Prior Art/Related Art\nThe ink jet recording method involves ejecting small droplets of recording liquid and landing them on a recording material such as paper (hereinafter referred to as a recording medium) to record information.\nThis recording method, because of its advantages such as low noise, low running cost, small size of apparatus and ease with which color printing is implemented, has found wide use in printers and copying machines.\nThis kind of ink jet recording generally employs ink whose major component is water that contains water-soluble, high boiling point solvent, such as glycol, for the purpose of preventing drying and clogging. Recording the ordinary paper using such ink, however, may result in poor fixing of ink on the paper or non uniform image presumably due to uneven distribution of loading materials and sizes applied to the surface of the paper. When a color image is formed, in particular, ink droplets of two or more different colors are successively ejected overlapping the previous ones before the droplets that have landed become fixed to the recording medium, so that colors may get blurred at boundary portions in the image between different colors or mixed unevenly causing feathering. As a result, a satisfactory printed image may not be obtained.\nTo solve the above-mentioned problem, a method has been proposed which involves purchasing a particular recording medium applied over its entire surface to a specified thickness with a substance that fixes a dye in the recording liquid and ejecting the recording liquid against the recording medium.\nWith such a method, however, satisfactory recording cannot be realized unless the above-mentioned special recording mediums are purchased.\nAnother method is known which sprays a colorless treatment liquid, that renders the dye in the recording liquid insoluble, over a wide area on one surface of the recording medium irrespective of the positions where the recording liquid droplets are landed.\nBecause the treatment liquid is also applied to portions not related to ink ejection, the amount of the treatment liquid is far greater than the amount of ink used, increasing the running cost and the size of the apparatus.\nA first object of this invention is to provide an ink jet recording apparatus as well as an ink cartridge and an ink jet head unit mountable on the apparatus, which can solve the above-mentioned conventional problems, can perform water-resistant, anti-bleed processing on a recorded image swiftly and accurately with a minimum amount of processing liquid, and can perform this processing without requiring a recording head of special construction for ejecting the processing liquid.\nA second object of this invention is to provide an ink jet recording apparatus as well as an ink cartridge and an ink jet head unit mountable on the apparatus, which can realize, without degrading operability, high-quality, high-reliability image recording that exhibits improved water resistance for ordinary paper and produces no feathering or color bleed at the boundaries between different colors in color recording.\nIn order to achieve these objects, in a first aspect of the present invention, a recording apparatus, which performs recording on a recording medium by using a coloring material-containing ink accommodated in a first accommodating portion and processing liquid accommodated in a second accommodating portion to make insoluble or coagulate the coloring material in the ink, comprises:\na means to hold a first ink jet head for ejecting the ink and a second ink jet head for ejecting the processing liquid in line with each other and opposite the recording surface of the recording medium; and\na control means to control the ejecting of the processing liquid from the second ink jet head in connection with the ejecting of the ink from the first ink jet head.\nHere, the first ink jet head and the second ink jet head may be arranged in the order of the second ink jet head and the first ink jet head in the direction in which these two heads perform at least the image recording so that the processing liquid can be ejected from the second ink jet head prior to the ejecting of the ink from the first ink jet head.\nThe first ink jet head and the second ink jet head may be combined together to form an ink jet head unit.\nThe ink jet head unit may be removably mounted on a carriage that reciprocally moves in a direction in which the unit moves when performing the image recording.\nThe ink jet head unit mounted on the carriage may be removably provided with the first accommodating portion that accommodates the ink containing the coloring material and with the second accommodating portion that accommodates the processing liquid to make insoluble or coagulate the coloring material in the ink.\nIt may further comprise a processing liquid supply means to supply the processing liquid to the second accommodating portion.\nThe processing liquid supply means may be removably fixed on a guide along which the carriage reciprocally moves and, only when the carriage moves to a predetermined position, is connected to the second accommodating portion to supply the processing liquid thereto.\nThe processing liquid supply means may further include a pressurizing pump that uses a driving force of the carriage to supply the processing liquid to the second accommodating portion.\nThe pressurizing pump may include bellows that pressurizes the interior of the processing liquid supply means when the second accommodating portion is connected to the processing liquid supply means.\nThe processing liquid supply means may include a soaking body that retains waste liquid of the processing liquid or the ink.\nThe first ink jet head may have an electro-thermal conversion body as an element that generates thermal energy for ejecting the ink.\nIn order to achieve these objects, in a second aspect of this invention, an ink jet head unit, comprises:\na first ink jet head to eject a coloring material-containing ink; and\na second ink jet head to eject processing liquid to make insoluble or coagulate the coloring material in the ink;\nwherein the first ink jet head and the second ink jet head are integrally combined.\nHere, the second ink jet head and the first ink jet head may be arranged in this order in a direction in which the both heads move when performing at least image recording.\nIn order to achieve these objects, in a third aspect of this invention, an ink cartridge, comprises:\na first accommodating portion that is removably fixed to a first ink jet head which ejects a coloring material-containing ink, the first accommodating portion accommodating the ink to be supplied to the first ink jet head.\nHere, the first accommodating portion may be divided into a first portion containing a absorbing body that absorbs and retains the ink and a second portion that accommodates the ink and communicates with the first portion through small holes.\nThe first accommodating portion may be divided according to the kind of the coloring material in the ink.\nThe first accommodating portion may be divided as claimed in the density of the ink.\nIn order to achieve these objects, in a fourth aspect of this invention, an ink cartridge, comprises:\na second accommodating portion that is removably fixed to a second ink jet head which ejects processing liquid to make insoluble or coagulate a coloring material contained in ink, the second accommodating portion accommodating the processing liquid to be supplied to the second ink jet head.\nHere, a surface tension of the processing liquid may be smaller than the surface tension of the ink.\nIn order to achieve these objects, in a fifth aspect of this invention, an ink cartridge, comprises:\na first accommodating portion that accommodates ink to be supplied to a first ink jet head which ejects a coloring material-containing ink; and\na second accommodating portion that accommodates processing liquid to be supplied to a second ink jet head which ejects the processing liquid to make insoluble or coagulate the coloring material in the ink;\nwherein the first accommodating portion and the second accommodating portion are integrally combined.\nThe ink may contain anionic dye.\nThe ink may contain at least anionic compound and pigment.\nThe processing liquid may contain a cation substance including a low-molecular component and a high-molecular component.\nWith this invention, because the recording head is so arranged as to allow the processing liquid to be sprayed prior to the ejecting of ink onto the recording medium, it is possible to apply a sufficient amount of processing liquid for recording to the recording medium. According to the recording data, the processing liquid is sprayed over only the locations on the recording medium where the ink is to be ejected, so that the ink lands on the applied processing liquid. This makes it possible to form a desired stable image by using only the minimum required amount of processing liquid, thereby optimizing the amount of processing liquid applied prior to the ink ejecting and lowering the running cost and the manufacture cost of the apparatus.\nFurther, because the recording heads of the same construction and configuration are employed, there is no need to change the design of the recording head according to whether the droplets to eject are the processing liquid or the ink, thus improving the development efficiency and productivity.\nMoreover, the recording head for ejecting the ink containing a coloring material and the recording head for ejecting the processing liquid for coagulating the coloring material in the ink may be integrally held by a frame, or the nozzle for the ink containing the coloring material and the nozzle for the processing liquid that coagulates the coloring material in the ink may be provided in the same recording head. This construction makes it possible to eject the processing liquid and the ink immediately following it in each line during each scan of the recording heads. This in turn allows the original purpose of the processing liquid to be accomplished even when the image data transfer from the host computer is halted.\nFurthermore, the recording heads for ejecting the processing liquid may be arranged on both sides of the recording head for ejecting the ink. This construction makes it possible to apply the processing liquid to the recording medium even during the bi-directional scan of the carriage, which in turn allows the output time to be reduced by the bi-directional recording and the recording speed to be improved.\nIn the ink jet recording apparatus of this invention, the colorless or light-colored transparent liquid (referred to as a processing liquid) containing a compound to render dye in the ink insoluble is sprayed onto the recording medium according to the image information, followed by ink of one color or two or more colors being ejected against the recording medium according to the image information. The apparatus has a dedicated tank for the processing liquid whose capacity is equal to or greater than that of the colored ink tank and which is installed in the apparatus body. The apparatus is so constructed as to make easily recognizable to an operator the necessity for replacing the tank of the processing liquid. High reliability of the recording apparatus is thus realized.\nThese and other objects, effects, features and advantages of the present invention will become more apparent from the following description of embodiments thereof taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates to filter designs, and more particularly to a variable passband autoregressive moving average (ARMA) filter using a non-causal filter design.\nA problem common to several signal processing applications, including radio frequency (RF) tuning, audio tone controls, variable transient filtering, sampling rate conversion, jitter/wander compensation and/or measurement and the like, is that a zero-pole based digital filter, such as an infinite impulse response (IIR) and finite impulse response (FIR) combination, needs to have a variable frequency response at a given sample rate, or vice versa—constant frequency response at a variable sampling rate. In video there is such a need for a variable bandwidth video luminance filter in order to reduce sensitivity of transients seen by a signal out-of-range detector/alarm. Such a filter needs to have a nominal lowpass response, as given by IEEE-P205, and via one control parameter continuous range from full bandwidth to some small fraction of nominal bandwidth. An additional requirement is that the filter needs to have greater computational efficiency—faster/less expensive—than existing filter methods in order to have realtime performance in a video waveform monitor, for example.\nPrior digital variable bandwidth filters include the use of:\n1) either only FIR or only IIR filters;\n2) analog filter simulation with variable C, L, R, gyrator, etc.;\n3) discrete bandwidth selection of ARMA design.\nThe first class satisfies the continuously variable bandwidth requirement, but generally requires more computation to approximate a pole with many zeroes and vice versa. The second class also satisfies the continuously variable bandwidth requirement, but generally does not work well when the cut-off frequency approaches Nyquist—at which point instability may take place for higher order filters or extra processing may be required to prevent it. Also in the second class mapping the passband control parameter to new filter component values is not always readily apparent or may be difficult or impractical to implement. The third class does not satisfy the continuously variable bandwidth requirement.\nWhat is desired is a stable, continuously variable bandwidth controllable via one parameter, digital filter for processing signals from full bandwidth to a small portion of the bandwidth."}
{"text": "1. Field of the Invention\nThis invention relates to wire spring forming machines.\n2. Description of the Prior Art\nKnown in the prior art as relating to the structure herein are the following U.S. Patents:\na. U.S. Pat. No. 3,351,101 to Halverson et al requires the use of an arbor to coil a spring and can coil only in a left hand or right hand direction. The spring coiled is wound about an arbor from which the spring is required to be unwound in a reverse coiling action for release from the arbor.\nThe machine herein feeds wire continuously in a wire feeding direction, the spring coil being formed free of being wound and retained about an arbor or without the use of an arbor and each spring upon being fully formed is cut off from the wire feed stock without any reversal action. During the time that Halverson is releasing a spring from an arbor by a reverse coiling action, the machine herein forms another spring.\nb. U.S. Pat. No. 3,025,889 to Clay assigned to Baird Corporation. Here two feed rolls are used with a groove formed on each roll requiring a change in feed rolls to accommodate changes in wire size and in addition liner sections and deflecting tools have to be changed.\nThe machine herein is adapted to receive a range of wire sizes and the only change required is that of a simple replacement of a quill having the appropriate size of passage therethrough to align with a corresponding size passage in the feeding apparatus of the machine.\nc. U.S. Pat. Nos. 3,025,889, 3,025,890 and 3,025,891 to Clay assigned to Baird Corporation disclose the use of a segmental gear which provides limited forward feeding motion to the extent of the gear which is intended to feed the length of wire required for a given spring and the gear is then reversed to its starting position for the next wire feeding cycle and thus there is no continuous forward feeding action of the wire feeding elements as is the case with the machine described herein.\nDuring the interval of time in which the segmental gears in the above indicated patents are reversed to a starting position, the machine herein will produce another spring and thus have on the order of twice the production output in producing wire springs as with the case of the above described patents."}
{"text": "1. Field of the Invention\nThe present invention is directed to expandable abrasive/cutting heads for rotation atherectomy devices.\n2. Description of the Related Art\nAtherectomy is a non-surgical procedure to open blocked coronary arteries or vein grafts by using a device on the end of a catheter to cut or shave away atherosclerotic plaque (a deposit of fat and other substances that accumulate in the lining of the artery wall). For the purposes of this application, the term “abrading” is used to describe the grinding and/or scraping action of such an atherectomy head.\nAtherectomy is performed to restore the flow of oxygen-rich blood to the heart, to relieve chest pain, and to prevent heart attacks. It may be done on patients with chest pain who have not responded to other medical therapy and on certain of those who are candidates for balloon angioplasty (a surgical procedure in which a balloon catheter is used to flatten plaque against an artery wall) or coronary artery bypass graft surgery. It is sometimes performed to remove plaque that has built up after a coronary artery bypass graft surgery.\nAtherectomy uses a rotating shaver or other device placed on the end of a catheter to slice away or destroy plaque. At the beginning of the procedure, medications to control blood pressure, dilate the coronary arteries, and prevent blood clots are administered. The patient is awake but sedated. The catheter is inserted into an artery in the groin, leg, or arm, and threaded through the blood vessels into the blocked coronary artery. The cutting head is positioned against the plaque and activated, and the plaque is ground up or suctioned out.\nThe types of atherectomy are rotational, directional, and transluminal extraction. Rotational atherectomy uses a high speed rotating shaver to grind up plaque. Directional atherectomy was the first type approved, but is no longer commonly used; it scrapes plaque into an opening in one side of the catheter. Transluminal extraction coronary atherectomy uses a device that cuts plaque off vessel walls and vacuums it into a bottle. It is used to clear bypass grafts.\nPerformed in a cardiac catheterization lab, atherectomy is also called removal of plaque from the coronary arteries. It can be used instead of, or along with, balloon angioplasty. Atherectomy is successful about 95% of the time. Plaque forms again in 20-30% of patients.\nSeveral devices have been disclosed that perform rotational atherectomy. For instance, U.S. Pat. No. 5,360,432, issued on Nov. 1, 1994 to Leonid Shturman, and titled “Abrasive drive shaft device for directional rotational atherectomy” discloses an abrasive drive shaft atherectomy device for removing stenotic tissue from an artery, and is incorporated by reference herein in its entirety. The device includes a rotational atherectomy apparatus having a flexible, elongated drive shaft having a central lumen and a segment, near its distal end, coated with an abrasive material to define an abrasive segment. At sufficiently high rotational speeds, the abrasive segment expands radially, and can sweep out an abrading diameter that is larger than its rest diameter. In this manner, the atherectomy device may remove a blockage that is larger than the catheter itself. Use of an expandable head is an improvement over atherectomy devices that use non-expandable heads; such non-expandable devices typically require removal of particular blockages in stages, with each stage using a differently-sized head.\nIn the years since the '432 patent, there has been a significant effort to improve the expandable head, with many devices using centrifugal force or other devices to drive portions of the abrasive head radially outward at high rotational speeds.\nFor all of these devices, the abrasive head includes an abrasive that has a single set of properties. For instance, an abrasive burr may include abrasive particles of a particular size or a particular distribution of sizes. Or, a particular head may have a cutting effect on the blockage, rather than a grinding effect.\nThere may be some instances when a practitioner requires two different abrading heads for a single blockage. For instance, a particular blockage may have hard plaques, which may be effectively removed by sanding or scraping, as well as soft lesions, which may be effectively removed by slicing or cutting. The cutting head may have different properties than the scraping head.\nIf a practitioner wants to use a first abrasive, then use a second abrasive having different properties than the first abrasive, the practitioner must remove the device with the first abrasive, then insert the device with the second abrasive. This removal of one catheter and insertion of another catheter is time-consuming, inconvenient, expensive, and requires additional parts that must be manufactured, shipped, inventoried, and maintained with the atherectomy device.\nAccordingly, there exists a need for a rotational atherectomy abrading head that can exhibit two different sets of properties. Such a head would reduce the expense, time and burden of using additional heads for the rotational atherectomy device."}
{"text": "1. Field of the Invention\nThe present invention relates to the generation of color models which characterize display devices.\n2. Description of the Related Art\nTypically, display devices, such as liquid crystal displays (LCD) or cathode-ray tubes (CRT), produce colors on a screen by combining different amounts of red, blue and green light. The actual color produced from specific combinations of RGB light varies from display device to display device. Therefore, in order to consistently reproduce colors on different displays, manufacturers of display devices create color models that characterize the output of a particular display device in terms of a standard color coordinate system. For instance, Commission Internationale de l'Enclairage (CIE) XYZ color coordinates are often used.\nConventionally, a display device characterization is represented by the following equation:\n      [                            X                                      Y                                      Z                      ]    =            [      M      ]        ·          [                                                  R              ′                                                                          G              ′                                                                          B              ′                                          ]      R′G′B′ are radiometric scalar values that are related to device RGB values according to generally non-linear relationships commonly called “gamma curves.” The matrix M is a 3×3 tristimulus matrix which relates the linear R′G′B′ values to XYZ values.\nConventionally, the tristimulus matrix M is populated according to the following equation:\n  M  =      [                                        X                          r              ,              max                                                            X                          g              ,              max                                                            X                          b              ,              max                                                                        Y                          r              ,              max                                                            Y                          g              ,              max                                                            Y                          b              ,              max                                                                        Z                          r              ,              max                                                            Z                          g              ,              max                                                            Z                          b              ,              max                                            ]  Each X, Y, and Z value in matrix M corresponds to a measured XYZ value when the maximum R, G and B values are applied to the device. Once determined, the tristimulus matrix M maps any linearized R′G′B′ value used by the display device to the actual XYZ color that the display device is expected to reproduce.\nThis conventional model for characterizing display devices is built upon the assumption that the measurements are done in a dim environment where ambient lighting has little effect on the measurements. This assumption may not always prove true. For example, many display devices will produce a non-zero amount of light when displaying the color black. This black point measurement is often called flare. In addition, a non-zero flare value can be caused by higher than optimal ambient lighting. In either case, the flare skews the results of the measurements.\nAttempts to measure or estimate flare, and correct for its effect, are well-known in the art. However, these attempts have generally been unsatisfactory. For instance, the effect of flare has been built into a general gain-offset-gamma-offset (GOGO) model for gamma curves. This model attempts to correct for the effects of flare in gamma curves. However, the effect of this model is limited since each gamma curve can only act on each channel independently. Other efforts have been made to correct for flare in the tristimulus matrix, but these efforts still utilize the conventional formula for populating the matrix.\nAs such, the conventional method does not fully address problems associated with flare, since only measurements at maximum RGB values are used to populate the tristimulus matrix. Even if flare is accounted for, the measurements themselves may be inaccurate, and therefore, the resulting tristimulus matrix may not provide the most accurate device characterization."}
{"text": "A liquid dispensing apparatus for decorative or industrial purposes may be designed for dispensing liquid in a uniform sheet. Such a liquid dispensing apparatus may be used for a variety of applications because the apparatus can be sized to produce many sizes of liquid sheeting. In addition, the liquid dispensing apparatus may be inexpensive as it may be constructed from common elements and in few manufacturing steps."}
{"text": "Optical fiber cable development, wherein the cable is capable of multichannel transmission, has led to a number of different cable configurations, prominent among which are bonded arrays of fibers forming a planar ribbon, and stacks of such ribbons within a core tube or sheath. In a typical ribbon array, a plurality of fibers, e.g., twelve, are held in spaced position parallel to each other by a suitable matrix, a configuration which simplifies construction, installation, and maintenance by eliminating the need for handling individual fibers. Thus, the splicing and connecting of the individual fibers can be accomplished by splicing and connecting the much larger ribbons provided that the fiber positions in the ribbon are precisely fixed and maintained.\nIn U.S. Pat. No. 4,900,126 of Jackson, et al., there is shown a bonded optical fiber ribbon which typifies, in many respects, the present state of the art for fiber ribbon. The ribbon comprises a plurality of longitudinally extending coated optical fibers disposed in a parallel array. A matrix bonding material fills interstices between adjacent fibers and extends about the array to form substantially flat protective surfaces. The ribbon may contain a large number of fibers such as, for example, sixteen, or it may have more or less, depending upon the particular application.\nCabling of fiber ribbon is usually accomplished by stacking a number of ribbons on top of each other and enclosing them loosely in a tube configuration while imparting a helical twist to the stack such as is shown in \"A Modular Ribbon Design for Increased Packing Density of Fiber Optic Cables\" by K. W. Jackson, et al., International Wire & Cable Symposium Proceedings 1993, pp. 20-26. In the manufacture of such a loose tube stacked ribbon configuration, the finished cable is wound upon a reel for storage and shipment. Quality control usually dictates that the transmission characteristics of such a cable wound upon a reel be measured for the cable in its reeled or storage condition. It has been found that in a reeled cable configuration, i.e., the fibers in the ribbon stack, exhibit a higher attenuation than does the same cable when laid out in a straight line configuration. Thus, it sometimes happens that a wound or reeled cable is rejected because of too high attenuation, when, in actuality, the cable in use would have an attenuation that was well within acceptable limits. This artificial packaging induced loss becomes greatly exacerbated with higher fiber count cables that are larger in sheath diameter. Such an increase in attenuation is attributable, therefore, to the winding of the cable on a drum, and is not predictable. That is, the increase may, in some instances, be large, while in other cases it is small. Thus, correlation of the attenuation increase with the size of the storage drum and the physical characteristics of the cable is difficult and unreliable. As a consequence, one cannot derive a precise \"increased attenuation factor\". Short of measuring cable attenuation with the cable in a straight-line configuration, which is totally impractical for the lengths of cable involved, it is desirable that a means be found for minimizing the transient increase in attenuation of fibers in a typical reel configuration.\nA similar problem can arises when an optical fiber, buffered fiber, or ribbon, such as dispersion compensating fiber, is wound upon a drum or reel. Such fiber, wound upon a conventional reel, exhibits greater or higher loss than when the fiber is laid out straight. This too is a result of the compressive forces placed on the fiber in wound condition. In dispersion compensating fiber, the fiber usually remains upon the reel in use, hence, an extra component of lose is introduced into the active circuitry."}
{"text": "The present invention relates to a state-of-charge measuring method, and more particularly, to a state-of-charge measuring method using multilevel Peukert's equation for battery-powered products such as electric vehicles.\nAccurate battery state-of-charge information is very important for the driver of an electric vehicle since, by knowing the exact amount of energy available from the batteries, the driver of the vehicle can travel without the fear of running out of energy before reaching the destination.\nThe state-of-charge of a battery is a function of used capacity and available energy. The used capacity can be easily calculated by simply multiplying the average discharge current and total time. On the other hand, determining the available energy of a battery is not straightforward because the available energy depends on the magnitude of the average discharge current. Low average discharge current increases the available energy, and high average discharge current decreases the available energy. And therefore, to calculate the available energy accurately, a sophisticated algorithm such as Peukert's equation must be used. In other words, to measure the state-of-charge of a battery accurately, the usage pattern such as driving profile must be taken into account.\nThere exist many methods for measuring battery state-of-charge. They include open-circuit voltage check, measurement of electrolyte's specific gravity, measurement of current under load and current accumulation (ampere-hour or Ah) method to name a few. However, all of the above methods do not take into account the usage pattern, and therefore, are susceptible to a large degree of error.\nAs mentioned earlier, one method that takes into account the usage pattern in determining the state-of-charge of a battery is Peukert's equation. Peukert's equation is used in conjunction with Ah method to reduce error. By using Peukert's equation, the available energy of a battery can be calculated with good accuracy, and consequently, the state-of-charge can be determined more precisely.\nIn a conventional method using Peukert's equation, the usage range of a battery is established and the maximum (I.sub.H) and minimum currents (I.sub.L) are set according to the battery's usage range, and Peukert's constants (n and K) are obtained by using the maximum and minimum currents. Peukert's equation calculates the available capacity (Ah.sub.-- available) of a battery as follows: EQU Ah.sub.-- available=KI.sub.avg .sup.(1-n)\nwhere n and K are Peukert's constants and I.sub.avg is the average discharge current. The constants n and K are found as follows: ##EQU1## and EQU K=I.sub.H.sup.n t.sub.H =I.sub.L.sup.n t.sub.L\nwhere I.sub.H and I.sub.L represent upper and lower current values, respectively, of the current range in which a battery is normally operated and t.sub.H and t.sub.L are discharge times at I.sub.H and I.sub.L, respectively. Using Peukert's equation, the state-of-charge of a battery is calculated as illustrated in FIG. 1. Referring to FIG. 1, in step 101, discharge current I(t) is read. In step 102, the average discharge current (I.sub.avg ) is calculated by using equation (1). ##EQU2##\nIn step 103, the capacity drawn from the battery (Ah.sub.-- used) is calculated by using equation (2) and in step 104, the available capacity (Ah available) is calculated by using equation (3) EQU Ah.sub.-- used=I.sub.avg .times.t (2) EQU Ah.sub.-- available=KI.sub.avg .sup.(1-n) (3)\nFinally in step 105, the state-of-charge is calculated by using equation (4) and the process returns to step 101. ##EQU3##\nPeukert's equation is a preferred method for applications having widely fluctuating loads such as electric vehicles since it can determine the available capacity under any discharge conditions. However, since the conventional Peukert's equation uses a maximum current I.sub.H and a minimum current I.sub.L as only references, its accuracy deteriorates when the average discharge current is faraway from either I.sub.H or I.sub.L as shown in FIG. 3. Therefore, to improve the accuracy over the entire range, it is necessary to introduce another reference point, i.e., an intermediate I.sub.M, between the maximum current I.sub.H and the minimum current I.sub.L as shown in FIG. 3."}
{"text": "1. Field of The Invention\nThe present invention relates to a device for fixedly mounting a heat sink to a central processing unit (CPU) for dissipating heat generated therefrom.\n2. The Prior Art\nSince the processing speed of computer CPUs has increased resulting in higher temperatures, heat dissipation of the CPU becomes an increasingly important issue. To solve this problem, a number of devices for mounting heat sinks to CPUs have been proposed as in U.S. Pat. Nos. 4,745,456, 5,019,940 and 5,594,624 and Taiwan Patent Application Nos. 83217100-1, 83217100-2, 84207163, 84212747, 85200861 and 85206293. This prior art discloses a mounting device which can fasten a heat sink, a CPU and a CPU connector together, thereby enabling the heat sink to dissipate heat generated by the CPU. However, the prior art mounting devices have the following disadvantages.\n1. The mounting devices disclosed in U.S. Pat. Nos. 4,745,456 and 5,019,940 and Taiwan Patent Application Nos. 83217100-2, 84212747 and 85200861 consist of too many components and have a complicated structure.\n2. The prior art devices disclosed in U.S. Pat. No. 5,594,624 and Taiwan Patent Application Nos. 83217100-1, 84207163 and 85206293 each have a simpler structure, however, specially designed tools are required to fasten/detach the heat sinks to/from the CPU. Therefore, the cost of the heat sink mounting devices is increased and assembly thereof is further complicated.\nHence, an improved heat sink mounting device is needed to eliminate the above mentioned defects of current heat sink mounting devices."}
{"text": "The concept of mixing materials by utilizing \"static\" or motionless mixing has been known for sometime. In the past 4 years, two designs have been available on the market one of which consists of a series of helical elements in a tube or pipe, and the other of which utilizes a complex series of tubular channels. Both of these structures divide and recombine a stream in geometric progression, and within a relatively short distance, feed stock material is thoroughly and predictably mixed. (\"Automation\", February 1972 \"Motionless Mixers\"). The helical element type of device is clearly disclosed in U.S. Pat. No. 3,286,992 to Armeniades et al. dated Nov. 22, 1966. These devices are called motionless mixers because they have no moving parts. Relative movement of the fluid with respect to the motionless mixing elements is, however, achieved by the flow of fluid within the conduit.\nOther efforts at blending various materials include the patent to Heyl et al. No. 2,601,018 wherein the blending tube contains a single perforated sheet metal spiral member substantially throughout its length, the perforated surface of the spiral member extending from wall to wall of the blending tube. Rogers in U.S. Pat. No. 2,628,864 is disclosing an aerosol paint spraying device taught the use within the spray tube of a spiral form member formed either of twisted wire or a helical ribbon of metal. Andrews et al. in U.S. Pat. No. 2,710,250 taught the mixing of fluids with a series of orifice members in a conduit. Grubb et al. U.S. Pat. No. 2,863,649 taught an apparatus for mixing on a small scale of compositions having a short period of coexistence when mixed and utilizing a rotating mixing rod having a small wire spirally wound around it and including at regularly spaced intervals spiral notches. Another device is taught in the patent to Thomas et al. U.S. Pat. No. 3,089,683 wherein an elongated tubular body is provided with a series of diffusers which create a turbulent flow of the liquids thereby ensuring a complete mixture of the liquids prior to ejection through the outlet. U.S. Pat. No. 3,203,371 to Mosey teaches a machine for whipping of confectionary filling utilizing in the nozzle thereof a baffle which comprises a strip of chrome steel twisted into a helical form and having a plurality of transverse slits to provide a multplicity of teeth or tongues which extend more or less radially from the axis of the helical bent strip.\nThe present invention is distinguished from these prior art devices in that the motionless mixing element is a singular structure of far simpler geometric configuration then that heretofore proposed and therefore much less costly in either fabrication or disassembly and cleaning than prior art structures."}
{"text": "Surface mounted light fixtures typically come in two parts—a back plate that is coupled to a wall or a ceiling, and a cover that is coupled onto the back plate. During installation, it is convenient to have the two parts near each other to accomplish the electrical connections. Often, while the installer makes the electrical connections, the installer has to hold the cover in one hand which in turn leaves the installer with only one free hand to make electrical connections. Since the installer cannot use both his/her hands to make the electrical connection, the installer may find this approach very inconvenient and inefficient.\nMechanisms for placing the cover in proximity to the back plate during installation do exist. For example, a hinge may be provided between the two parts. However, adding hinges to the cover and back plate may be cost prohibitive. Another example includes providing a tether between the two parts. However, adding tethers to the cover and the back plate may also be cost prohibitive. Furthermore, tethers may not rigidly hold the cover in proximity to the back plate and the installer still may be required to support the cover with one hand. Yet another example includes providing snaps that hold the cover and back plate together during installation. However, the installation of snaps are limited to plastic parts that can utilize the snaps. Accordingly, snaps are not a feasible option for light fixtures made using other types of materials.\nIn view of the foregoing shortcomings, there is a need for technology that holds the two parts of the surface mounted light fixtures rigidly and in close proximity to each other during the installation of a surface mounted light fixture. This background information is provided to reveal information believed by the applicant to be of possible relevance to the present disclosure. No admission is necessarily intended, nor should be construed, that any of the preceding information constitutes prior art against the present disclosure."}
{"text": "This invention relates to a process for making human antibody-producing B lymphocytes in tissue culture. More specifically, this invention relates to the antigen specific induction and regulation of antibody synthesis in cultures of human peripheral blood mononuclear cells. Such antibody producing cells can be used to make hybridomas using cell fusion techniques for producing monoclonal antibodies for therapeutic and diagnostic purposes.\nThe knowledge of mechanisms that govern the humoral immune response has been greatly advanced by the culture technique, initially developed by Mishell and Dutton for murine spleen cells, which permitted the study of antibody production in virto, Mishell et al., J. Exp. Med. 126: 423, 1967. Application of this technique developed for murine cells to the study of human lymphocytes has been less successful. Human B lymphocytes are obtained from peripheral blood, tonsils, spleens, and lymph nodes. While tonsil and spleen cells have been shown to produce antibodies against certain antigens [sheep red blood cells (SRBC)] in vitro under conditions which are similar to the conditions commonly used in experiments with murine spleen cells, peripheral blood lymphocytes (PBL) did not produce antibodies when cultured and tested in the same way, Hoffman, et al., Nature 243: 408, 1973. As peripheral blood lymphocytes (PBL) represent the only source of lymphocytes which is readily accessible, progress in the study of the humoral immune response in man has been hampered by the lack of technique which permits in vitro sensitization of peripheral blood lymphocytes.\nIf one disregards the well-established assay that measured the polyclonal production of antibodies which does not depend on antigen, Fauci et al., J. Exp. Med. 144: 674, 1926, success has been reported only by two groups. In the first case, infection with the Epstein-Barr virus (EB virus) was required to enable PBL to produce antibodies in specific response to SRBC or horse red blood cells (HRBC), Luzatti et al., Nature 269: 419, 1977. In the second case, Dosch et al., J. Immunol. 118: 302, 1977, no additional stimulus was used, but the hemolytic plaques in the standard assay has been questioned.\nWith respect to the use of EB virus, although it stimulates response in Human B lymphocytes, it apparently operates through an extraordinary pathway which cannot be reproduced with more benign mitogens. Furthermore, the usefulness of antibody producing cell lines generated using EB virus, for treatment of human patients is questionable as one would always be in danger of the effect of such viruses. EB virus is known to be tumor producing and also to cause mononucleosis. For these same reasons, laboratory personnel do not want to handle them either.\nIt has been found that one can use natural pathways of antibody production to generate antibody producing cells from human B lymphocytes. This avoids the need to deal with such extraordinary pathway mitogens as the tumor producing EB virus, with their attendant dangers to personnel working with the virus and patients receiving a possibly contaminated or otherwise dangerous end product. The production of antibody producing cells from Human B lymphocytes (or B-cells) is based on the realization that it is helper cells (macrophages and helper T cells) that support the B cell responses and it is suppressor lymphocytes, present in the B lymphocytes source which suppress or interfere with Human B lymphocytes activation. The effect of the suppressor lymphocytes can be overcome either by their destruction or removal from the Human B lymphocytes or by use of a B cell mitogen to overcome their effect."}
{"text": "1. Field of the Invention\nThe present invention relates to a fuel injection system for an internal combustion engine, and more particularly, to a fuel injection system including a fuel pump which periodically transmits pressurized pulses of fuel from the fuel pump through a conduit to a fuel injection valve.\n2. Description of the Prior Art\nA wide variety of fuel injection systems include a cam driven fuel pump which periodically transmits pressurized pulses of fuel from the fuel pump through a conduit to one or more fuel injector valves. Typically these fuel injector valves are spring biased to a closed position and open when the pressure at the fuel injector valve exceeds a predetermined pressure level. As the amplitude of the pressure pulse present at the fuel injector valve decreases below a predetermined threshold level, the force provided by the spring bias closes the needle valve of the fuel injection valve and terminates the injection of fuel into an engine cylinder.\nA 1936 United Kingdom patent specification Ser. No. 443,124 discloses an invention by Amery which relates to an improvement for fuel injection systems of the type described above. In the system disclosed by Amery the fuel injector valve is opened when the fuel pressure at the valve reaches a predetermined level. An auxiliary cam-operated control valve is provided to generate a fluid pressure which is applied to the back side of the fuel injector valve to forcefully close the injector valve at predetermined times independent of the fuel pressure at the discharge end of the fuel injector valve. In accordance with the Amery system, commencement of fuel injection depends on the timing of the mechanically driven fuel pump while the termination of the fuel injection depends on pressure applied to the closing piston within the fuel injector valve under the control of the cam operated control valve.\nAn inherent characteristic of engine speed driven fuel injection systems is that the pressure available for injection decays rapidly at low engine speeds due to the slower pumping action of the fuel pump.\nU.S. Pat. No. 4,069,800 (Kanda et al) discloses a fuel injection system which utilizes accumulators to maintain a uniform fuel supply pressure at the fuel injector valve. A spring biased servo valve within the fuel injector valve controls the opening and closing of the fuel injector valve. U.S. Pat. No. 4,036,192 (Nakayama) discloses a related fuel injection system which is intended to eliminate dribbling of fuel from fuel injection nozzles following the termination of each fuel injection pulse.\nU.S. Pat. No. 3,952,711 (Kimberley et al) discloses a fuel injection nozzle with independent opening and closing control. In this common rail system which utilizes an accumulator to provide a uniform source of fuel pressure at the fuel injection valve, the opening of the nozzle valve is controlled by a compression spring in a conventional manner, but an independent means is provided to control the nozzle valve biasing forces during opening and closing of the nozzle. U.S. Pat. No. 3,416,506 (Steiger) discloses another type of common rail fuel injector system which utilizes a pressure accumulator to maintain a uniform pressure level at the fuel injector valve."}
{"text": "The conventional rear derailleur for shifting the driving chain has hitherto been well-known, which ordinarily comprises a fixing member, a linkage mechanism provided with a base member, two linkage members and a movable member, and a chain guide provided with a guide pulley, a tension pulley and a pulley plate provided with an urging plate which extends toward the path of the chain travelling from the guide pulley to each sprocket at a multistage sprocket assembly and is opposite to the chain. The linkage mechanism, which fixes the base member to the fixing member, is mainly swingable by a first horizontal shaft, interposes between the linkage mechanism and the fixing member a first spring which biases the linkage mechanism to move away from the center of the multistage sprocket assembly, pivotally supports the chain guide to the movable member in relation of being swingable through a second horizontal shaft, and interposes between the chain guide and the movable member a second spring which applies tension to the driving chain, so that the first and second springs balance in tension with each other to control the position of the guide pulley with respect to each sprocket at the multistage sprocket assembly.\nThe rear derailleur employing the aforesaid first and second springs is disclosed in, for example, the Japanese Patent Laid-Open Gazette No. Sho 58-149,882, which is so constructed that the first and second springs apply the predetermined tension to the chain and balance in tension with each other to set the chain guide in position, in detail, the guide pulley at the chain guide in position with respect to each sprocket at the multistage rear sprocket assembly. Also, the chain is shifted from a smaller diameter sprocket to a larger diameter one at the multistage sprocket assembly through the urging surface at the pulley plate coming into contact with the chain to forcibly urge the chain toward the larger diameter rear sprocket. In addition, generally an operating wire is fixed to one linkage member and an outer sheath for guiding the operating wire is supported to the base member.\nIn the driving stage that the chain is stretched between the smaller diameter sprockets at the multistage rear sprocket assembly and at a multistage front sprocket assembly, when the chain is shifted from the smaller diameter front sprocket to a larger diameter front sprocket by operating a front derailleur, a linkage mechanism L, as schematically shown in FIG. 9, swings counterclockwise around a first horizontal shaft P.sub.1 as shown by the arrow in FIG. 9 and along the plane perpendicular to a hub shaft and the chain guide CG swings in the same manner as the linkage mechanism, both the linkage mechanism L and chain guide CG shifting from the positions shown by the dotted lines to that shown by the dot-and-dash lines respectively. In this case, the outer sheath O supporting position on the linkage mechanism is displaced from the point X.sub.1 to X.sub.2 in FIG. 9 as the linkage mechanism L swings, so that a curvature of a curving portion of the outer sheath O at the multistage rear sprocket assembly side is reduced from that shown by the solid line to that shown by the dot-and-dash line. Accompanied with this, the tension of operating wire increases to act on the linkage mechanism L to an extent of the tension increase, whereby the linkage mechanism L, even if the operating lever is not operated, deforms automatically axially inwardly of the multistage rear sprocket assembly.\nThe conventional rear derailleur provided with the first and secnod horizontal shafts parallel to each other is constructed to make the path of swinging motion of the linkage mechanism parallel to that of the chain guide, so that when the chain is shifted from the smaller diameter front sprocket to the larger diameter one, the guide pulley GP at the chain guide CG shifts radially along the plane perpendicular to the hub shaft HS at the multistage rear sprocket assembly S as the linkage mechanism L and chain guide CG swing around the first and second horizontal shafts P.sub.1 and P.sub.2 respectively. The displacement increases the tension of the operating wire so as to deform the linkage mechanism L axially inwardly of the multistage rear sprocket assembly S, whereby the guide pulley GP is displaced from the position shown by the solid line to that shown by the dot-and-dash line in FIG. 15. In other words, the guide pulley GP is displaced toward a second rear sprocket S.sub.2 of a larger diameter and adjacent to a first rear sprocket S.sub.1 carrying the chain. As a result, the chain guided by the guide pulley GP to the first rear sprocket S.sub.1 is slanted between the guide pulley GP and the first rear sprocket S.sub.1 to come into contact with the urging surface at a pulley plate of the chain guide CG to generate a sound. In order to prevent the soundmaking, there is the trouble that an operating lever for the rear derailleur must intentionally be restored. Especially, when a positoning mechanism is provided at the rear derailleur or the operating lever therefor, the conventional derailleur has a large defect."}
{"text": "A lightning protection system protects a structure from damage due to lightning strikes, either by safely conducting the strike to the ground, or preventing the structure from being struck. Lightning protection systems are commonly installed on structures, buildings, trees, monuments, bridges and water vessels to protect from lightning or electrical discharge damage. Lightning protection systems are also used to protect appliances, machinery, electrical systems and electronic equipment from lightning or electrical discharge damage. Most lightning protection systems are composed of a network of lightning rods and/or lightning arrestors (surge protectors), metallic cable conductors, and ground electrodes designed to provide a low impedance path for the lightning to travel through to the ground.\nA lightning rod is a metal strip or rod with a conductive cable or other low-resistance path to ground. A lightning rod, which is usually installed on a rooftop, provides a point well above the structure to be protected with a very good, earthed connection. Because of its position, shape, and conductivity, the lightning rod may draw energy (current) from a lightning discharge and diverts the energy to ground via the conductive cable to ground, thus preventing damage to the structure. FIG. 1, which is adapted from Tesla U.S. Pat. No. 1,266,175, shows an early lightning rod.\nA lightning arrestor is a device, which is typically used for protecting electronic or electrical equipment from lightning by diverting any surges of high-voltage electricity caused by atmospheric discharges to ground. The lightning arrestor includes an “active” element, which switches from a non-conductive state to a conductive state in response to a surge in voltage. The active element in its conductive state provides a short i.e. a path for the high voltage to go to ground, bypassing the electronic or electrical equipment to be protected. In other words, the active element acts as an over-voltage release valve. The active element of a lightning arrestor may, for example, be a metal oxide varistor, a transient suppression diode, a gas discharge tube, a spark gap, a crowbar (circuit) using a Zener diode driving the gate of a silicon-controlled rectifier (SCR) latch, or any other suitable device. For convenience in description, both lightning rods and lightning arrestors may be referred to as lightning arrestors herein.\nConsideration is now being given to improving traditional lightning protection systems."}
{"text": "1. Field of the Invention\nThe present general inventive concept relates to a portable information processing apparatus, and more particularly, to a portable information processing apparatus having an improved structure of an input unit and an output unit on a front panel.\n2. Description of the Related Art\nIn general, a portable information processing apparatus is used in processing various information and provides convenience to a user in his/her daily life. As an example of the portable information processing apparatus, there are an MP3 player, a personal digital assistant (PDA), a cellular phone, an electronic frame and a portable multimedia player (PMP), etc. Hereinafter, an ultra mobile personal computer (UMPC) will be described as a small, portable and multifunctional computer, i.e., the portable information processing apparatus.\nThe UMPC is a data processor, and has functions of a notebook, a tablet personal computer, a navigator, a portable storage, a portable multimedia player (PMP), a personal digital assistant (PDA), a digital multimedia broadcasting (DMB) player, etc. Such an UMPC is lightweight, small-sized and portable, so that it is convenient for a user to transport. The UMPC allows a user to input and output information or the like in various environments.\nTherefore, if an input unit is efficiently arranged in the UMPC, it will be convenient for a user to input and output information, data, etc."}
{"text": "1. Field of the Invention\nThe invention relates generally to the field of electronic circuits and more particularly to designs for SRAM memory cells that provide improved stability in comparison to conventional designs.\n2. Related Art\nComputer systems and other devices typically need to have means for storing information.\nThese means may include persistent storage devices for large amounts of data, as well as smaller memory systems for storing data that the computer or other device is currently using. The memory systems for storing currently used data include both read-only memory (ROM) and random access memory (RAM.)\nRAM is typically used as the working memory of a device. RAM is used by devices to store data that needs to be accessible by a processor, and also needs to be modifiable. That is, the data can be changed. By contrast, data stored in a ROM cannot be modified, but can only be read. There is a great demand for RAM in computers and other electronic devices because the more RAM a device has, the more data can be readily accessible to the device's processor. For example, in a computer, the availability of more RAM enables the computer to execute more (or larger) software applications without having to swap data between RAM and a persistent storage device, such as a hard disc drive.\nThere are various different types of RAM. For example, dynamic RAM, or DRAM, has often been used in computers. The “dynamic” aspect of DRAM refers to the fact that DRAM memory cells need to be periodically refreshed in order to maintain the data which is stored in the cells. If the DRAM cells are not refreshed, the data will be lost. Static RAM, or SRAM, is another type of memory that is often used in computers. That “static” aspect of SRAM refers to the fact that SRAM cells do not have to be refreshed in the same manner as DRAM cells.\nSRAM memory has a number of advantages over DRAM memory. As noted above, SRAM cells do not have to be refreshed in order to maintain the data that is stored in them. Additionally, SRAM is typically much faster than DRAM. For example, typical SRAM cells may have access times of about less than 1 nanosecond, while DRAM cells may have access times closer to 60 nanoseconds. Further, SRAM memories do not require pauses between accesses, so the cycle time to access SRAM cells is typically much shorter than the cycle time for accessing DRAM cells.\nAlthough it has a number of advantages, SRAM memory also has some disadvantages. For instance, SRAM is typically much more expensive to manufacture than DRAM. Because the cost of SRAM memory is much higher than DRAM, it is common for SRAM memory to be used as a memory cache, while DRAM is used for a processor's main memory.\nSRAM memory cells may also be unstable. That is, the data in the cells may actually be corrupted when the cells are read. This problem arises from the fact that SRAM cells are read by coupling the cells to pre-charged bit lines and allowing the cells to pull down these bit lines. In other words, a higher voltage on the bit line is coupled to a lower voltage in the SRAM cell, causing the bit line voltage to drop and the SRAM cell voltage to rise. The voltage drop on the bit line is detected and amplified to provide the data for the associated processor.\nThe voltage rise in the SRAM cell, however, may corrupt the data stored in the cell. (because the initial, low voltage corresponded to the “0” stored in the cell, and the higher voltage resulting from the access may cause the data to be ambiguous, or even to flip-flop, so that it is now a “1.”)\nIf the transistors that make up the components of the memory cell were exactly identical to each other, the probability that the data in the cell would be corrected in this manner would be relatively low. As a practical matter, however, the transistors are not identical, but instead have slight variations that cause them to have slight variations in their respective responses. For example, each transistor has an associated threshold voltage that affects the response of the transistor. Because of manufacturing differences between the transistors, there is a variation in these threshold voltages that, in turn, results in a variation in the transistor's responses. These variations are becoming increasingly significant because the variation in threshold voltages increases as the size of the transistors decreases. Thus, as the memory cells grow smaller, they are more susceptible to the instability problem.\nThe instability of SRAM cells is obviously problematic. It would clearly be desirable to provide a design for an SRAM cell that is more stable than conventional designs and is therefore less likely to be corrupted when the cell is read."}
{"text": "1. Field of the Invention\nThis invention pertains to casting tools and the techniques for forming an exact replica of an eroded area in an underground gasoline storage tank and repair of the eroded area.\n2. Description of Prior Art\nA wide variety of approaches have been attempted to investigate this erosion problem. One of the approaches has been sonic graphing of eroded areas, inspection by closed circuit TV cameras. Dentists employ related techniques in producing mouth impressions of damaged or missing teeth. A known, related method for determining wear or erosion in the mechanical fields pertains to the determination of defects in a spur gear employing a casting procedure, such as U.S. Patent to Sawyer, No. 2,601,703, entitled \"Method for Testing Surface Defects.\""}
{"text": "High dependability and reliability is a goal of all computer and software systems. Complex systems in general cannot attain high dependability without addressing crucial remaining open issues of software dependability. The need for ultra-high dependability systems increases continually, along with a corresponding increasing need to ensure correctness in system development. Correctness exists where the implemented system is equivalent to the requirements, and where equivalence can be mathematically proven.\nThe development of a system may begin with the development of a requirements specification, such as a formal specification or an informal specification. A formal specification might be encoded in a high-level language, whereas requirements in the form of an informal specification can be expressed in restricted natural language, “if-then” rules, graphical notations, English language, programming language representations, flowcharts, scenarios, goal-directed requirement documents, or even using semi-formal notations such as unified modeling language (UML) use cases.\nNatural language scenarios can be constructed in terms of individual scenarios written in a structured natural language. Different scenarios can be written by different stakeholders of the system, corresponding to the different views the stakeholders have of how the system will perform, including alternative views corresponding to higher or lower levels of abstraction. Natural language scenarios can be generated by a user with or without mechanical or computer aid. The set of natural language scenarios provides the descriptions of actions that occur as the software executes. Some of these actions will be explicit and required, while others can be due to errors arising, or as a result of adapting to changing conditions as the system executes.\nFor example, if the system involves commanding space satellites, scenarios for that system can include sending commands to the satellites and processing data received in response to the commands. Natural language scenarios might be specific to the technology or application domain to which the natural language scenarios can be applied. A fully automated general purpose approach covering all domains is technically prohibitive to implement in a way that is both complete and consistent. To ensure consistency, the domain of application might be specific-purpose. For example, scenarios for satellite systems might not be applicable as policies for systems that manufacture agricultural chemicals.\nAfter completion of an informal specification that represents domain knowledge, the system is developed. A formal specification is not necessarily used by the developer in the development of a system.\nIn the development of some systems, computer readable code may be generated. The generated code is typically encoded in a computer language, such as a high-level computer language. Examples of the languages include Java, C, C Language Integrated Production System (CLIPS), and Prolog.\nOne step in creating a system with high dependability and reliability can be verification and validation that the executable system accurately reflects the requirements. Validation of the generated code is sometimes performed through the use of a domain simulator, a very elaborate and costly approach that is computationally intensive. This process of validation via simulation rarely results in an unambiguous result and rarely results in uncontested results among systems analysts. In some examples, a system is validated through parallel mode, shadow mode operations with a human operated system. This approach can be very expensive and exhibit severely limited effectiveness. In some complex systems, this approach leaves vast parts of possible execution paths forever unexplored and unverified.\nDuring the life cycle of a system, requirements typically evolve. Manual change to the system creates a risk of introducing new errors and necessitates retesting and revalidation, which can greatly increase the cost of the system. Often, needed changes are not made due to the cost of verifying/validating consequential changes in the rest of the system. Sometimes, changes can be simply made in the code and not reflected in the specification or design, due to the cost or due to the fact that those who generated the original specification or design are no longer available.\nProcedures, considered as the essential steps or actions to achieve a result, can be used for the assembly of materials in factories, for servicing of spacecraft (whether by astronauts, robots, or a combination), for business operation, and for experiments in a laboratory, to name but a few. Procedures can be very complex, involving many interactions, may involve many actions happening in parallel, and may be subject to significant constraints such as the ordering in which activities must happen, the availability of resources, etc. In many complex procedures, human error commonly results in the entire procedure needing to be repeated ab initio. In some cases, such as servicing a spacecraft, recovery from some of the more serious errors that may occur may not be possible.\nAs a rapidly growing field, autonomic systems (autonomic computing and autonomic communications) is a promising new approach for developing large-scale complex distributed' computer-based systems. In autonomic computing, the needs of large scale systems management has been likened to that of the human autonomic nervous system (ANS). The ANS, through the self-regulation, is able to effectively monitor, control and regulate the human body without the need for conscious thought. The self-regulation and separation of concerns provides human beings with the ability to concentrate on high level objectives without having to micro-manage the specific details involved.\nThe vision and metaphor of autonomic computing is to apply the same principles of self-regulation and complexity-hiding to the design of computer-based systems, in the hope that eventually computer systems can achieve the same level of self-regulation as the human ANS. The majority of conventional systems address the “how” of autonomic systems involving the low-level internal implementation, such as defining autonomic managers that together with the component that is to be managed make up an autonomic element to exist in a collaborative autonomic environment to provide self-management of the system. However, these efforts do not directly address the high-level requirements of the systems that drive autonomic systems.\nFor the reasons stated above, and for other reasons stated below which will become apparent to those skilled in the art upon reading and understanding the present specification, there is a need in the art to produce a system that is a provably correct implementation of agent-oriented specification, with or without autonomic properties, and that assures the consistency of such agent-oriented specifications throughout the implementation, that precludes major discrepancies, and that does not require use of a theorem-prover to establish correctness of the implementation. There is a further need for a convenient way of generating a new system when an agent-oriented specification changes. There is also a need for an automated, mathematics-based process for validation of agent-oriented specification that does not require large computational facilities."}
{"text": "Such articles typically have some portion, usually the backing layer, liner, or baffle constructed of a liquid repellent film material. This repellent material is appropriately constructed to minimize or prevent the exudation of the absorbed liquid from the article and to obtain greater utilization of the absorbent capacity of the product. The liquid repellent film commonly used includes plastic materials such as polyethylene films and the like.\nAlthough such products are relatively inexpensive, sanitary and easy to use, disposal of a soiled product is not without its problems. With greater interest being placed in protecting the environment today, there is a need to develop materials that are more compatible with the existing and developing waste disposal technologies while still delivering performance consumers have come to expect. An ideal disposal alternative would be to use municipal sewage treatment and private residential septic systems. Products suited for disposal in sewage systems can be flushed down a convenient toilet and are termed \"flushable.\" While flushing such articles would be convenient, the liquid repellent material which normally does not disintegrate in water tends to plug toilets and sewer pipes. It therefore becomes necessary, although undesirable, to separate the barrier film material from the absorbent article prior to flushing.\nIn addition to the article itself, typically the packaging in which the disposable article is distributed is also made from a water resistant material. Water resistivity is necessary to prevent the degradation of the packaging from environmental conditions and to protect the disposable articles therein. Although this packaging may be safely stored with other refuse for commercial disposal, and especially in the case of individual packaging of the products, it is often more convenient to dispose of the packaging in the toilet with the discarded disposable article. However, in the cases where such packaging is composed of a water resistant material, plugging of the drains to the toilet typically results.\nDesirably, a commercial, flushable product should be relatively responsive to water and be transportable in a sewer system. Commercially available water-soluble polymers, such as polyethylene oxide (PEO), polyvinyl alcohol (PVOH), acrylamide polymers, acrylic acid-based polymers, and cellulose derivatives, possess the desired characteristics for flushability, such as water solubility and/or water dispersibility. However, due to their in-use degradability and storage degradation, these materials function poorly as components in personal care products. Other disadvantages are that these polymers are difficult to process and are substantially more expensive than polyolefins.\nThe requirements for a functional and flushable product provide a substantial challenge in finding suitable materials with the desired properties. In an attempt to overcome the flushability problem of a water resistant film the prior art has modified the water resistant polymer. One of the more useful ways of modifying polymers involves blending them with other polymers of different structures and properties.\nPolymer blends of polyolefins and poly(ethylene oxide) have been shown to be water modifiable at expectedly low weight % polyolefin levels. Such blends would be anticipated to be flushable when exposed to water in a toilet but do not possess the dry mechanical properties required for functionality in use. Moreover, the high content of poly(ethylene oxide) makes such materials prohibitively expensive for use in a disposable personal hygiene article such as a sanitary napkin, diaper and the like. Polymer blends of polyolefins and poly(ethylene oxide) containing greater than about weight percent of polyolefin are generally water resistant and are not water modifiable.\nIn view of the problems of the prior art, it remains highly desirable to provide a water modifiable film having a substantial portion of thereof composed of a polyolefin. More desirably, the water modifiable film should have greater than about 55 weight percent of a polyolefin. When dry, the film should have the mechanical properties necessary for functionality. When wet, the films should lose at least a portion of its mechanical properties which would render the film flushable and transportable in a sewer system. Such films could be used for making flushable barrier films for personal care products.\nIt is therefore an object of the invention to provide a polyolefin-containing film that is water modifiable or water-degradable which contains higher levels of polyolefin content. More specifically, it is an object of the invention to provide a polyolefin-containing film having greater than about 55 weight percent of a polyolefin and less than about 45 weight percent of poly(ethylene oxide) that is water-modifiable or water degradable."}
{"text": "Electrical loads for vehicles continue to escalate. At the same time, the overall package size available for the electrical generator continues to shrink. Consequently there is a need for a higher power density system and method of generating on-board electricity.\nThe increase in the demand for current to be produced by an alternator leads to a need for corresponding adaptation of the system for current regulation, and particularly of a system whereby the three-phase alternating current produced by an alternator is rectified, into a direct current, which can be stored in a battery of a vehicle or be used directly by the electrical circuit of the vehicle which is supplied with a direct current (DC) voltage.\nRectification of the three-phase alternating current is generally obtained by means of a rectifier bridge having six power diodes. Three of these diodes are the positive diodes, and are connected between the phase terminals of the stator windings of the alternator and the positive terminal B+ of the alternator which is connected to the battery and the electrical circuit of the vehicle. Three further diodes, namely the negative diodes, are connected between electrical ground or earth of the vehicle and the aforementioned phase terminals of the stator windings.\nThe diodes constitute the rectifier bridge and are subjected to high current. Hence, it is necessary to cool them in the most effective way possible. To this end, it is known to arrange the diodes on metal plates, which are arranged on the outside of the alternator and which constitute a dissipator for the heat produced by the diodes. The diodes are grouped on two carrier plates, one of which is reserved for the positive diodes, and the other for the negative diodes.\nThe rectifier diodes are connected to respective carrier plates, and these carrier plates are used as heat sinks for these diodes as well. The rectifier diodes are inserted by pressure in receiving bore holes of the carrier plate or heat sink, or are soldered to the carrier plate using appropriate solder alloys. The end wires connected to the rectifier diodes enable the rectifier diodes to be connected to external sources.\nHowever, under certain particularly severe operating conditions, it has become apparent that this cooling of the diodes, and in general terms the cooling of the whole of the current regulating means, can be insufficient to ensure long-term reliability of the alternator.\nThe heat sinks are typically constructed in the shape of a circle or crescent and are fastened in the same plane to the alternating current generator.\nIt is important that the bridge rectifiers must not only be able to withstand normal battery charging current, but must also be able to supply high electrical load currents demanded by the vehicle and do so at high, elevated ambient temperatures. These demanding situations may occur, for example, when the vehicle is stuck in stop and go traffic on a hot night. Under such conditions the electrical loads are high due to headlights, air-conditioning, engine cooling fan, brake lights and various other miscellaneous electrical loads. The underhood temperatures are also high due to the outside ambient air temperature and the stop and go driving schedule. Bridge rectifiers, as discussed, are limited in their ability to reliably function under such high current and high temperature conditions. In order for bridge rectifiers to handle these types of excessive currents and heat, it becomes necessary to utilize a bridge rectifier which has higher current handling capability. Due to the space limitations of the alternating current generator, it then becomes very difficult to provide such a bridge rectifier from a feasibility standpoint as well as at an economical cost.\nIncreasing the current capacity and heat dissipating characteristics of the bridge rectifier has included mounting of semiconductor diode chips onto first and second metallic heat sinks which are electrically insulated from each other by a thin sheet of electrical insulating material. The diode chips are then covered by a protective insulating coating after connection to the respective heat sink. One of the metallic heat sinks includes a finned area which is subjected to cooling air when the bridge rectifier is mounted to the generator. The heat sink with the plurality of fins includes twelve air passages. This type of bridge rectifier is shown in U.S. Pat. No. 4,606,000 to Steele et al., incorporated herein by reference.\nWith time and vibration the electrical integrity of the electronic circuitry can fail due to fatigue from mechanical stress. This produces a loss of electrical output from the machine. Further, several parts are required to provide the output connection and therefore the piece or part cost and assembly cost is high.\nThere is a need to improve the performance characteristics of prior art bridge rectifiers. In addition, there is a need for a carrier plate which increases the dissipation of heat from the diode and more efficiently cool the diode by facilitating increased surface area of the heat sink available to a diode resulting in an increased rate of dissipation of heat from the diode through the carrier plate in which the diode is disposed. Furthermore, a more robust package assembly capable of withstanding vibration while providing an effective thermal dissipation that is also cost effective is accordingly desired."}
{"text": "Draught beer is often preferred by consumers to bottled or canned beer. Draught beer is generally served at the counter of a public house out of a refrigerated keg provided with a fluid connection to a source of pressurized gas for driving the dispensing of the beer through a dispensing line fluidly connecting the keg to a dispensing tap, comprising a valve for controlling the flow out of said tap. In case of a temporary social event outside a public house, such as an outdoor event, wedding party, fair and the like, consumers like to be offered draught beer for consumption. Furthermore, above a critical volume of consumption, serving bottled or canned beer would be too expensive and would generate too much waste. For these reasons, roving or mobile beverage dispensing units, offering the same quality of beer as a draught beer served at a public house, were developed and brought to the market. They are designed to accommodate a keg or container containing the beer, with a source of pressurized gas, such as a pressure gas bottle or a compressor. The containers used can be traditional metal kegs as used in public houses, possibly but not necessarily of smaller dimensions, or can include so called bag-in-containers as disclosed e.g., in EP2146832, EP2148770, EP2148771, EP2152494 and the like.\nFor example, US2004/0226967 proposes a dispensing unit comprising a cooling chamber suitable for accommodating and cooling a beer keg, a hollow column supported on said cooling chamber and a dispensing head comprising a tap valve. A source of pressurized gas, such as a compressor or a CO2 cartridge is provided for ensuring the necessary pressure for driving the beer flow out of the keg. A dispensing tube fluidly connects the keg to the tap valve. For reasons of hygiene, the dispensing tube is disposable and must be changed with each new keg. In one embodiment, the dispensing line is even permanently coupled to the keg to ensure that it will not be used a second time. Upon use, a new keg can be installed into the cooling chamber, and fluidly connected to a source of pressurized gas, generally located in the same chamber. The dispensing line is either permanently coupled to the keg or must be coupled thereto, before it is run through a channel defined in the hollow column until the dispensing tube outlet reaches the dispensing head of the column and is engaged into the tap valve mechanism. This “bottom-up” insertion system, wherein the dispensing tube is installed starting from the keg (located at the bottom) all the way up to the dispensing head (located at the top) requires that the dispensing line be provided with a shut-off valve to prevent the flow of beer out of the keg before the dispensing line is in place in the tap valve. It is clear that providing a shut-off valve to a disposable tube increases substantially the cost of use of the system. Furthermore, it can be quite cumbersome to drive up a flexible dispensing line through the hollow column which outlet to the cooling chamber is positioned at the back thereof and can easily be appreciated when looking e.g., at FIG. 2 of US2004/0226967.\nIn order to facilitate the engagement of the dispensing tube into the tap valve, a rather critical operation which is difficult to control from the interior of the cooling chamber, WO2009/115928 suggests to allow the opening of the dispensing head so that the dispensing tube outlet emerging from the opening at the top of the column can be handled from outside the cooling chamber and engaged more comfortably into the tap valve mechanism.\nEP1982952 extends the idea of allowing the opening of the column to the entire length thereof. This solution greatly simplifies the “bottom-up” installation of the dispensing tube since it needs only be passed from the interior to the exterior of the cooling chamber through a short channel crossing the top board of the cooling chamber before it can be handled from outside the cooling chamber, instead of having to drive it from the inside of the cooling chamber all the way up to the dispensing head.\nIn spite of the various solutions proposed to simplify it, the “bottom-up” installation of a disposable dispensing tube remains cumbersome since the user must crouch and engage the head and shoulders into the cooling chamber to access the opening connecting the cooling chamber to the dispensing column inner channel, push up the flexible tube either all the way up to the dispensing head like in US2004/0226967, or only until the outlet of the tube reaches the opening in the column as in WO2009/115928 and in EP1982952, at which point it must be grabbed from the outside before it falls back all the way down into the cooling chamber. Since this operation must be repeated with each new keg installed into the cooling chamber, if the installation of the tube is too uncomfortable, users may become reluctant to use such roving beverage dispensing unit.\nThe change of dispensing tube with every new keg remains a rather delicate operation, often to be performed in the dark, in noisy and crowded environment, and also in stressful conditions. The present invention proposes a solution to greatly simplify the loading of a new dispensing tube into a tapping column to bring the dispensing unit into operational condition whenever a new keg is loaded."}
{"text": "1. Field of the Invention\nThe present invention relates to a program mode acess control system by which a plurality of input output devices are controlled by one input output control device and a rapid program mode access is made between this input output device and a central processing device.\n2. Description of the Related Art\nThe construction of an input output control device which controls an input output device such as a magnetic tape device, a magnetic disc device, a magnetic floppy device, a line printer, etc., usually comprises a microprocessor, and corresponds to the construction of the input output device. Therefore, a standardization of the input output device is desired.\nThe conventional processing device, for example, has the construction, described as follows. A processor bus couples a central processing unit, a main memory device and a channel device, a plurality of input output control devices are connected via the channel device and a common bus, and each input output control device is connected to an input output device such as a magnetic tape device, a magnetic disc device, a magnetic floppy device, a line printer, or a display device etc.\nEach input output control device is formed by an internal control portion formed by a microprocessor, etc., a receiver driver, a control register, an input output interface circuit, etc. One or a plurality of the same kind of input output devices are connected to the input output interface circuit.\nAs described above, the conventional input output control device provides a control register formed by a plurality of registers, sets the control information in the control register and reads out the content thereof to carry out the desired control operation. Thus, the control system due to the wired logic is conventionally used. However, due to recent developments in microprocessors, the control system can now be operated with firmware, and a flexible control for the input output device is now possible, which leads to a reduction of the hardware needed.\nThe construction of the control system on the firmware by such a wired logic corresponds to that of the input output device, such as the magnetic tape device, magnetic disc device, display device, etc. Thus when a plurality of the input output devices are controlled, the control system is restricted to that of one kind of input output device.\nHowever, control of a plurality of the input output device should be made by the same input output control device. In such case, the operating speed of the microprocessor (the internal control portion) in the input output control device cannot follow the access operation for the high speed control register due to command checks from the central processing device side. Further, when using hard logic, much hardware is required for the control register. Thus a unified control of a plurality of different kinds of input output devices becomes difficult."}
{"text": "1. Technical Field\nThe present invention relates to semiconductor circuits in general, and in particular to domino type logic circuits. Still more particularly, the present invention relates to a domino logic circuit having a clocked precharge.\n2. Description of the Prior Art\nA domino type of logic circuit simplifies digital logic by connecting a number of transistors together in series to implement digital combination logic. For example, a domino logic circuit implements a logic AND function by cascading a P-channel transistor with several N-channel input transistors in series. During operation, the P-channel transistor is clocked to precharge an output node of the circuit to a predetermined logic state. Depending on the logic state at the inputs of the N-channel input transistors, the output node either remains at its pre-charged state or is pulled low through the series of N-channel input transistors by a clocked N-channel transistor connected to ground.\nIn accordance with the logic AND function, if all of the N-channel input transistors are driven by a logic high signal, the output node of the domino logic circuit will be a logic low. Conversely, if any one of the N-channel input transistors is driven by a logic low signal, the output node of the domino logic circuit will remain at its precharged logic high state. Because no inversion function is performed with this conventional arrangement, an inverter is generally employed at the output of the domino logic circuit to perform a logic inversion function such that an overall logic AND function is realized.\nThe present disclosure describes an improved domino logic circuit that is able to provide an output signal faster than the conventional domino logic circuit described above."}
{"text": "1. Field of the Invention\nThe invention is in the technical field of crop protection products which can be employed against harmful plants, for example in crop plants, and which comprise, as active compounds, a combination of at least two specific herbicides.\n2. Description of Related Art\nThe document EP-A-131258 discloses sulfonylureas and their salts and also their use as herbicides and/or plant growth regulators.\nThe efficacy of these herbicides against harmful plants in the crop plants is at a high level, but depends in general on the application rate, the formulation in question, the harmful plants or spectrum of harmful plants to be controlled in each case, the climatic conditions, the soil conditions and the like. Another criterion is the duration of action, or the breakdown rate of the herbicide. If appropriate, changes in the sensitivity of harmful plants, which may occur upon prolonged use of the herbicides or within geographic limitations, must also be taken into consideration. The compensation of losses in action in the case of individual harmful plants by increasing the application rates of the herbicides is only possible to a certain degree, for example because such a procedure frequently reduces the selectivity of the herbicides or because the action is not improved, even when applying higher rates. In some cases, the selectivity in crops can be improved by adding safeners. In general, however, there remains a need for methods to achieve the herbicidal action with a lower application rate of active compounds. Not only does a lower application rate reduce the amount of an active compound required for application, but, as a rule, it also reduces the amount of formulation auxiliaries required. It both reduces the economic input and improves the ecological compatibility of the herbicide treatment.\nOne possibility of improving the application profile of a herbicide can consist in combining the active compound with one or more other active compounds. However, the combined use of a plurality of active compounds frequently causes phenomena of physical and biological incompatibility, for example a lack of stability in a coformulation, decomposition of an active compound, or antagonism of the active compounds. What is desired are, in contrast, combinations of active compounds having an advantageous activity profile, high stability and, if possible, a synergistically improved action, which allows the application rate to be reduced in comparison with the individual application of the active compounds to be combined."}
{"text": "The present invention relates to a MEMS component, in the layer structure of which at least one sound-pressure-sensitive diaphragm element is formed, which spans an opening or cavity in the layer structure. The diaphragm element is structured and the deflections of the diaphragm element are detected with the aid of at least one piezosensitive circuit element, which is situated in the area of the attachment of the diaphragm element to the layer structure.\nA particular advantage of piezosensitive MEMS microphone components is that—in contrast to capacitive MEMS microphone components—they may be equipped very easily with a “wake-up” functionality. They may thus be provided very easily so that they only consume power if needed, i.e., for example, only if a specific sound level is exceeded. The power consumption of piezosensitive MEMS microphone components in the “always-on” operating mode is thus significantly less than that of capacitive MEMS microphone components.\nPiezosensitive MEMS microphone components are described in U.S. Patent Application Pub. No. 2014/0084395 A1. These MEMS microphone components each include a structured sound-pressure-sensitive diaphragm element having a piezoelectric layer structure for signal detection. The structuring of the diaphragm element is used here for pressure equalization between the two sides of the diaphragm element and is to promote the sound-pressure-related diaphragm movement. For this purpose, a uniform grid arrangement of passage openings in the entire diaphragm area is proposed, or also trench slits, by which the diaphragm element is divided into paddle-like segments. The piezoelectric layer structure is situated in each case in the edge area of the diaphragm element, where the greatest elongation or the greatest mechanical stresses occur when the diaphragm element is deflected."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical interrogation/reading system and method for detecting biological substances using one or more grating-based waveguide sensors (e.g., resonant waveguide gratings) that in one embodiment can be incorporated within a microplate.\n2. Description of Related Art\nToday studies associated with detecting a biological substance (e.g., cell, drug, chemical compound) using resonant waveguide gratings (RWGs) are fast becoming the technology of choice in academia and industry. In such studies, a reading system is used to couple light into the RWG and also used to analyze the light that is reflected from the RWG in order to determine whether or not a biological substance is present on the surface of the RWG. To detect the biological substance, the reading system analyzes the reflected light to locate its resonant wavelength/angle which corresponds to a certain refractive index that is indicative of whether or not the biological substance is located on the surface of the grating-based waveguide sensor. Unfortunately, the reading system used in industry today suffers from a major drawback wherein it is difficult for a user to tune the RWG and reading system so that the reading system can properly interface with the RWG. In particular, it is difficult for a user to tune the RWG so that its resonant wavelength/angle is at a desired spectral location so that the resonant wavelength/angle can be easily detected by the reading system. This drawback is especially troublesome when a multi-channel reading system is used to interface in a multiplexed manner with an array of RWGs. Accordingly, there is a need for a way to make it easier to tune a reading system and RWG to address the aforementioned shortcoming and other shortcomings in the prior art. This need and other needs are satisfied by the reading system, RWG and method of the present invention.\nThe present invention includes a reading system and at least one grating-based waveguide sensor that interact with one another to enable the detection of biological substances (e.g., cells, drugs, chemical compounds). In one embodiment, the reading system includes a light source (e.g., laser, diode) for directing a light beam into the grating-based waveguide sensor and a detector (e.g., spectrometer, CCD imaging device) for receiving a reflected light beam from the grating-based waveguide sensor and analyzing the reflected light beam so as to detect a resonant wavelength/angle which corresponds to a predetermined refractive index that indicates whether a biological substance is located on the surface of the grating-based waveguide sensor. The grating-based waveguide sensor is tuned to have a resonant wavelength/angle at a predetermined spectral location by adjusting a skew angle defined as an angle between a plane of incidence of the light beam directed into the grating-based waveguide sensor and a grating vector which is perpendicular to the lines of a diffraction grating within the grating-based waveguide sensor. In another embodiment, the reading system is capable of performing a multiplexed interrogation of an array of grating-based waveguide sensors."}
{"text": "The present invention relates generally to an electronically commutated direct-current motor (ECM).\nMotors of this kind are used, inter alia, to drive miniature fans (cf. EP-A1-0 908 630 and corresponding U.S. Pat. No. 6,013,966, FEHRENBACHER et al). For various reasons, it may be desirable to operate such a motor without a Hall generator, and for that purpose to commutate it with the voltage that is induced, during operation, by the rotor in the stator winding. This is known, for example, from U.S. Pat. No. 4,156,168, VOGEL, but in the case of this known motor the direction in which it will start after being switched on is not certain, and this known motor is therefore suitable only for specific applications in which rotation direction is not important.\nOne object or the invention is therefore to provide new electronically commutated motor whose startup rotation direction is certain. In accordance with the invention, the motor used is one having a rotor that, when the motor is currentless, has at least one mechanical preferred position. The current in the motor is controlled using a bistable multivibrator that has, at startup, an electrical preferred position which is adapted to the mechanical preferred position of the rotor. At startup, the result of the electrical preferred position is that the stator is excited in such a way that the rotor starts, from its mechanical preferred position, in the correct rotation direction. Because a separate rotor position sensor is eliminated, a motor of this kind has a simple configuration with good efficiency, since the power consumption for a rotor position sensor is eliminated. The invention is therefore particularly advantageous for miniature and subminiature motors in which the power consumption of a rotor position sensor, for example a Hall IC, would greatly reduce the electrical efficiency; and it is highly suitable, for example, for motors in which the electronic components are arranged separately from the actual motor (i.e. stator with stator winding, and rotor).\nFurther details and advantageous developments of the invention are evident from the exemplary embodiment described below and shown in the drawings, which is in no way to be understood as a limitation of the invention."}
{"text": "Hydroelectric power generation, wind-power generation, photovoltaic power generation and the like, which can be used to attempt to reduce carbon dioxide or improve other environmental problems by using inexhaustible natural energy, have received much attention. Among these, photovoltaic power generation has seen a remarkable improvement in performance such as the power generation efficiency of solar cell modules, and an ongoing decrease in the price, and national and local governments have worked on projects to promote the installation of residential photovoltaic power generation systems. Thus, in recent years, spread of photovoltaic power generation systems has advanced considerably.\nPhotovoltaic power generation directly converts solar energy to electric energy using a silicon cell semiconductor (solar cell element). The solar cell element as used herein undergoes a decrease in function when directly brought into contact with air from the outside. Therefore, a solar cell element is placed between sealing materials so as to provide a buffer as well as to prevent incorporation of foreign materials or infiltration of moisture. This sealing material is required to have various performance criteria, such as that the material is transparent and does not inhibit the power generation by light (transparency), that the material does not elute or disintegrate under the action of heat (heat resistance), that the material has good adherence to a protecting material such as glass or a back sheet (adhesiveness), and that the material does not undergo significant deterioration or yellowing under sunlight (durability). Thus, in order to satisfy these requirements, various formulations of blending in various components have been investigated. For example, there is known, as a representative formulation with consideration for transparency, heat resistance, adhesiveness, flexibility, moldability, durability and the like, a mixture of an ethylene-vinyl acetate copolymer combined with a peroxide and a silane coupling agent is used as a solar cell sealing material (see, for example, Patent Document 1).\nSuch a solar cell sealing material needs to employ a two-step process including a step of producing a sheet of an ethylene-vinyl acetate copolymer with various additives incorporated, and a step of sealing a solar cell element using the sheet thus obtained. In this step of producing a sheet, since it is necessary to perform molding at a low temperature at which the organic peroxide included for crosslinking does not undergo decomposition, the extrusion molding speed cannot be increased. Furthermore, in the step of sealing a solar cell element, it is necessary to perform an adhesion process which takes time and has two steps, including a step of performing temporary bonding in a laminator over several minutes to several tens of minutes, and a step of performing main bonding in an oven at a high temperature at which the organic peroxide is decomposed for several tens of minutes to one hour. Therefore, the production of a solar cell module requires effort and time, and also is one contributory factor for increasing the production cost.\nIn light of such circumstances, investigations are being undertaken on substitute materials which do not require the use of an organic peroxide such as described above, can significantly improve the production efficiency for solar cell modules, and have excellent characteristics as sealing materials for solar cells. Specifically, there has been suggested a solar cell element sealing material using an ethylene-unsaturated carboxylic acid copolymer having an unsaturated carboxylic acid content of 4 mass % or more and a melting point of 85° C. or higher, or an ionomer thereof (see, for example, Patent Document 2).\nFurthermore, a press sheet has been disclosed, which is obtained by performing press molding using, as a sealing material that constitutes a solar cell module, an ethylene copolymer composition containing an ethylene-methacrylic acid copolymer and an amino group-containing trimethoxysilane at a predetermined ratio (see, for example, Patent Documents 3 and 4).\nOn the other hand, a solar cell uses various backing films on the side opposite to side of the solar cell element where sunlight is incident, and the amount of addition of a silane coupling agent in the sealing material may be increased in order to adjust the adhesiveness to the backing films\nPatent Document 1: WO 2006/095762 A1\nPatent Document 2: Japanese Patent Application Laid-Open (JP-A) No. 2000-186114\nPatent Document 3: JP-A No. 2004-31445\nPatent Document 4: JP-A No. 2001-144313"}
{"text": "It known to provide planters with lamps to illuminate plants and some planters also provide some illumination for the surrounding area.\nHowever, none of the prior planters mount a lamp adjacent a corner formed by the junction of the horizontal surface or tread of a step on which the planter stands and the vertical riser of the next step so that the light beam preferentially illuminates that corner."}
{"text": "Aircraft tractors or tugs are used in combination with a tow bar attached to the nose wheel of an aircraft for transporting the aircraft on the ground, such as ramp and taxi positions. The tug typically has a hitch at the rear having upper and lower apertured plates. The tow bar frontally terminates with an apertured tow eyelet. For connecting, the tug operator maneuvers the hitch into proximity with the tow bar, while the eyelet is aligned with the hitch by an assistant on the ground. Inasmuch as the tug structure prevents direct visual observation, the tug operator is reliant on the assistant for movement instructions. When properly aligned the assistant inserts a latch pin through the plates and the eyelet to complete the coupling. The operation requires accurate and timely communication between the parties to accomplish a time efficient and safe connection. These requirements are often compromised by operating conditions resulting in a large incidence of personnel injuries. "}
{"text": "This invention generally relates to crucible covers for multiple pocket electron beam sources used in vacuum coating systems in which individual crucible pockets containing coating materials are heated to produce coating vapors. The invention further relates to cooperating cover and crucible combinations that isolate heated pockets from all other pockets by forming a line-of-sight blocking baffle.\nVacuum coating or deposition of materials is a process commonly used in the manufacture of a variety of products including, but not limited to, semiconductors, optics, and the formation of thermal barriers on high temperature components. One method for coating materials having very small vapor pressures, such as metals or ceramics, is by placing the coating material and a workpiece to be coated in a vacuum chamber, reducing the pressure (typically in the range of 10xe2x88x927 to 10xe2x88x924 Torr), and heating the material. The resulting coating vapor streams outwards from the material and coats the surface of the workpiece. Due to the low vacuum chamber pressure, vaporized coating materials travel nearly unimpeded, following an approximately line-of-sight trajectory. The most common sample heating method uses directed electron beam energy.\nMany manufacturing steps involve depositing multilayered coatings of a multitude of materials. Generally it is important to be able to control the coating purity and deposition thickness to achieve desired results. In addition, many of the benefits of producing multiple coating layers on a workpiece are achieved when the coating steps are carried out sequentially under vacuum, and the trend in coating technology has been towards obtaining purer, more uniform and controllable coating thickness of multiple materials. An apparatus for carrying out sequential coating is disclosed in U.S. Pat. No. 2,482,329, which describes a mechanism for sequentially moving crucible pockets in line with a single heating source in a single vacuum chamber. While the mechanism of that patent was used in conjunction with electric resistance heating, the same mechanism works just as well with other sources, such as electron beams (see, for example, U.S. Pat. No. 4,632,059 to Flatscher et al.). The use of electron beams, which can be positioned using controllable magnetic fields, allows for the option of keeping the samples in place while deflecting the heating source, as disclosed in U.S. Pat. No. 5,792,521.\nOne major problem with the open material containers disclosed in the previously referenced patents is cross-contamination between crucibles. Although the majority of coating materials is directed away from its source, some material invariably diffuses back towards the source, contaminating the other source materials. This reduces the purity of the coating material and limits the usefulness of multiple coating chambers. One solution to this contamination problem is to incorporate a noncontacting crucible cover or movable lid with an opening that allows vapor from the heated sample to coat the workpiece, while shielding all of the other samples contamination. A typical cover arrangement is shown in U.S. Pat. No. 4,748,935 to Wegmann, which discloses an electron beam heated vapor source with such a cover. The arrangement is that of a flat, fixed cover and a multiple pocket crucible that rotates the samples through a heating region, allowing each pocket in turn to become a heated pocket. The cover and crucible are separated by a small gap that limits, but does not eliminate, line-of-sight contamination. In this configuration coating deposits that build up on the lid, interfering with the desired operation of the cover and requiring much maintenance. In addition, this type of cover can result in a new path for contaminationxe2x80x94cover deposits can be scraped off or flake off into crucibles containing different materials. Nearly all covers known in the art rely on close tolerances to reduce, but not eliminate, line-of-sight contamination, and thus they still have the disadvantage of providing a path for contamination.\nSome of the limitations of noncontacting covers were addressed in U.S. Pat. No. 4,944,245 to Stoessl et al. In Stoessl et al., a cover contacts the crucible surface in the vicinity of the heated sample to hinder the flow of material from the heated sample to any unheated samples. The cover and crucible are connected so that crucible rotation lifts the cover, allowing for sample selection. Although this approach reduces contamination by vapors migrating from heated to unheated samples, it is well known by those in the art that coating material rapidly accumulates on cover and crucible surfaces located within line-of-sight and particularly on those surfaces that are also near heated samples. In addition, the lifting mechanism described by Stoessl et al. occupies a large fraction of the crucible surface area, limiting the number of samples which can be accommodated. The coupled rotation and lifting mechanism of Stoessl et al. also allows for multiple contacts per material selection if the rotation is not between adjacent crucibles, which may increase the opportunity for cross-contamination.\nIn both the Wegmann and Stoessl et al. configurations, deposits have a tendency to rapidly accumulate thick layers of coating material that can easily flake or be scraped off and contaminate adjoining samples during rotation of the crucible. In addition, the close tolerances between cover and crucible in both of those references can be violated by deposits, requiring equipment maintenance limited by the ability of the cover mechanism to operate properly. It is desired to have an improved crucible cover that can accommodate deposits while reducing contamination.\nThe prior art solution to reducing cross-contamination between pockets in multiple pocket sources relied on either closely spaced flat covers or covers that sealed with the crucible. Those solutions rely on maintaining close tolerances in regions where tolerances are difficult to maintain due to rapid build up of unwanted deposits. The present invention circumvents the problems in prior art systems by relying on a noncontacting baffle arrangement that is unaffected by unwanted deposits to a much greater degree. The present invention comprises a cover and crucible with matching features that, in combination, form a baffle to block the line-of-sight between heated and unheated crucible pockets, thereby reducing contamination. The rotation of the multiple crucible and lifting of the cover are controlled by separate mechanisms to provide flexibility in the design and operation of the source. Selection of any of the multiple pockets is achieved by rotating the crucible to the pocket of interest with the cover clear of the crucible. Cross-contamination protection is achieved by lowering the cover over the crucible to form a line-of-sight baffle after the selected crucible is properly aligned. In a first embodiment, the portion of the baffle nearest the heated pocket has a trough or grooves in the crucible with a cover lip that protrudes into the groove. This configuration allows the normal build up of coating materials on the cover top surface without causing degradation in performance from contamination or interference of moving parts. In a second embodiment there is a barrier that extends upwards from each of crucible pockets, while the cover protrudes below the barrier level, blocking line-of-sight between heated and unheated crucibles. In addition, the covers of the present invention lift away from the crucible prior to crucible rotation. In contrast to the prior art cam-based crucible lid combinations of Stoessl et al., in which the cover is lifted as the result of crucible rotation, the present invention reduces the amount of crucible surface area needed to implement the cover arrangement by providing a mechanism to lift the cover that does not depend on the rotation of the crucible.\nThe noncontact baffle arrangement of the present invention is an improvement over prior art covers are configured to either a) limit the line-of-sight between heated and unheated pockets with a close tolerance between cover and crucible, or b) forming a seal between cover and crucible. Both of those prior art covers have problems resulting from the build up of coating materials on the cover that tend to reduce their performance or that require frequent maintenance. In addition, some prior art covers place the edge of the cover near the top of the crucible pocket, thus allowing for the possibility of contamination due to deposits dropping directly into crucible pockets.\nBy providing a light tight, line-of-sight blocking baffle with clearance to allow for the build up of coating material, the present invention circumvents many of the problems with the prior art, allowing the cover to provide more contamination reduction for longer periods of time.\nIt is an advantage of the invention that a crucible cover for a multiple pocket evaporative source can be operated to greatly reduce cross-contamination of coating materials and that materials coating back onto the source not reduce the efficacy of the cover to reduce contamination or become a source for contamination.\nIt is yet another advantage of the invention that a multiple pocket crucible forms a line-of-sight blocking baffle in cooperation with a movable crucible cover, reducing cross-contamination between pockets without crucible-to-cover contact between heated and unheated crucible pockets.\nIt is an advantage of the invention that contamination is reduced with a crucible cover that can be lifted away from a crucible prior to rotation.\nIt is an advantage of the invention that it can be used for coating of materials that cold weld on contact by providing a contamination barrier that will not allow deposits to weld the crucible to the cover.\nIt is another advantage of the invention that materials that coat the cover do not hinder the movement of a turreted crucible.\nIt is a further advantage of the invention that the cover is lifted away from the crucible by an external lifting mechanism.\nIt is yet a further advantage of the invention that improved reliability and length of time between servicing will reduce the amount of scrap material and reduce the cost of operating coating equipment.\nIt is an advantage of the invention that it is simple and inexpensive to manufacture and service and that it is easy to use.\nAll patents cited herein are hereby incorporated by reference in their entirety for all purposes. Additional objects, advantages, aspects and features of the present invention will become apparent from the description of preferred embodiments, set forth below, which should be taken in conjunction with the accompanying drawings, a brief description of which follows, and in part will become apparent to those skilled in the art upon examination of the following or may be learned by practice of the invention."}
{"text": "This invention relates generally to filters, and more specifically to a filter bottom construction employed to support filter media in liquid filters.\nMulti-block filter bottoms of the type disclosed in U.S. Pat. No. 3,110,667, issued to Stuppy, have been employed with success in liquid filter systems, e.g., water and wastewater filter systems. The blocks are disposed next to each other in parallel adjacent rows, and then are cemented or grouted in place to form a support structure for gravel and/or other finely divided filter media.\nThe upper walls of the blocks are provided with passages through which liquid can pass during both the filtering and backwashing operations. The interior of each block is divided into upper and lower laterals separated by a horizontal partition, but interconnected by a plurality of vertical ports formed through the partition.\nThe above-described multi-block construction is commonly employed in a filter structure having a central flume extending below a concrete bottom wall upon which the multi-block construction is secured. The blocks that are vertically aligned with the flume have at least a portion of their bottom walls removed to provide openings into the lower laterals. These lower laterals communicate with the flume through connecting ports formed through the concrete supporting structure.\nDuring conventional filtration the influent moves in a downstream direction through the filter media, and then through the passages formed in the top walls of the blocks forming the filter bottom. The liquid then moves from the upper laterals into the lower laterals throught the vertical ports in the horizontal partitions, and then along the lower laterals to the cut out regions in the blocks overlying the central flume. At this location the effluent (i.e., the filtered liquid) passes into the flume through the connecting ports in the concrete supporting structure, and then out of the filter.\nDuring backwashing the backwash liquid is introduced into the flume and passes through the connecting ports in the concrete supporting structure into the cut out regions in the blocks overlying the flume. At this point the liquid moves into, and is distributed laterally along the lower laterals. The distributed backwash liquid passes through the ports in the horizontal partitions separating the upper and lower laterals in the various blocks, and then through the passages in the top walls of the blocks and through the filter media to provide the backwashing function.\nAlthough the above-described prior art system provides excellent liquid distribution during backwashing, while permitting the desired flow of influent during filtration, a problem does exist in the repair and maintenance of the system in the event that one or more of the blocks become clogged. For example, if the liquid distribution ports formed in the horizontal partition between the upper and lower laterals in one or more of the blocks become blocked, or if the laterals themselves become blocked, the distribution of backwash liquid can be adversely affected. This can cause uneven flow of backwashing liquid through the filter media, resulting in an undesired disruption of the media bed. Unfortunately, it is not always easy to determine exactly where the blockage exists merely by visually observing the movement of the backwash liquid. In fact, in attempting to locate a blockage it is not uncommon to actually break out one or more blocks that are perfectly good, on the mistaken belief that those blocks actually were the ones that were malfunctioning.\nThus, although the multi-block filter bottom constructions of the prior art are well suited for use during filtering and backwashing operations, they do not permit the easy, non-destructive detection of blockages therein."}
{"text": "According to the known art, magnetic steel sheets having a layer that improves the electrical insulation are used, for example, in electric drives for the design of stators. The materials used are regulated by the standard EN 10106 (1995). The materials named in this standard give a wide-ranging product range in order that the demands of different applications can be satisfied. The usable materials range from low-alloyed steel, with outstanding magnetic permeability, good thermal conductivity and good stamping properties, to higher-alloyed steels having very low remagnetization losses even at higher frequencies. As alloying constituents, the alloys in the standard contain copper (<=0.02%), manganese (<=1.2%), silicon (0.1-4.4%), aluminum (0.1-4.4%), the sum formed from the silicon content and twice the aluminum content being <5%, phosphorus (<=0.15%), tin (<=0.2%) and antimony (<=0.2%). Iron forms the basis of the alloy.\nCoatings which improve the insulation between the individual steel sheet layers and the processability have been developed for improving the properties of the magnetic steel sheets. The specific properties of the material used have to take into consideration influencing variables, such as corrosion protection, electrical insulation, influence on the stamping properties, heat resistance or weldability. Coatings for magnetic steel sheets can be gathered from the standard EN 10342 (2005).\nThe magnetic steel sheets available in the aforementioned standards, and the coatings thereof, cannot withstand all fields of use, however, as has been shown. Particularly when the magnetic steel sheets are exposed to highly corrosive media, e.g. sour gas (high hydrogen sulfide content), these magnetic steel sheets are at great risk of corrosion."}
{"text": "An antenna is a well known arrangement for radiating or receiving electromagnetic waves. While antennas are available in a variety of shapes and sizes, they all function based on the same basic principles. In the reception mode, an antenna intercepts a propagating electromagnetic wave, which then induces an electronic signal within the antenna. The electronic signal can be then fed into an integrated circuit that deciphers the signal. In the transmission mode, an antenna receives an electronic signal through a feed line, which then induces a field surrounding the antenna that results in the formation of a free-space propagating electromagnetic wave. The antenna's features such as its dimensions can be obtained by reference to its operation frequency, radiation patterns, loss, gain, and the like. Antennas are typically made from metallic materials and have a wide variety of configurations. One known configuration is a dipole antenna that includes two conductive bodies of equal length each receiving an input signal at one end thereof. The two conductive bodies of a typical dipole antenna are elongated bodies that are aligned with one another. Each body may be one-quarter of the wavelength of the target wavelength which is to be transmitted or received by the antenna.\nAntennas are prevalently used in wireless devices such as cell phones and the like to direct incoming and outgoing electromagnetic waves between a free space and a transmission line. Antennas are also used in radio frequency identification device (RFID) applications.\nAn RFID device that includes an antenna is usually referred to as an inlay. An inlay may include an antenna as well as a transponder, which is an integrated circuit for deciphering signals sent to the inlay and received by the antenna and also for sending a signal to the antenna which is then transmitted by the antenna. The inlay antenna may be tuned (i.e. sized) to communicate at a certain target frequency with a transceiver which is sometimes referred to as the interrogator. The interrogator typically includes an antenna for communication with the RFID inlay. An inlay may be active or passive. An active inlay would include its own power source such as a battery, while a passive inlay would receive its power from an external source such as an interrogator.\nA magnetic coupler that employs a terminated transmission line can be used in encoding of RFID-enabled labels, tickets, tags, cards or other media. U.S. Pat. Nos. 7,425,887 and 7,190,270 disclose RFID printers/encoders which employ single transmission line couplers for communication with RFID inlays.    PTL 1: U.S. Pat. No. 7,425,887 U.S. Pat. No. 7,190,270 U.S. Pat. No. 7,348,885"}
{"text": "Conventional repair and manufacturing facilities typically comprises numerous supports for holding individual sections or modules of gas turbine engines during assembly and disassembly. In this way the weight of the engine may be distributed over different points of support along the length of the engine. The supports are typically configured to externally hold the various engine sections in appropriate relative positions to permit assembly thereof. Each engine module is mounted to its dedicated support which is time consuming. Multiple jack stands need to be moved during assembly and weight of each engine module is typically supported by overhead cranes during most of the operations. The engines are typically support in a fixed position rendering access to higher or lower components difficult. The supports along the length of the engines also limit engine accessibility for the assembler/worker.\nThere is thus a need for a new ergonomic engine mounting arrangement to permit assembly and maintenance operations on gas turbine engines in a cost effective and time efficient manner."}
{"text": "1--Field of the Invention\nThe present invention is generally concerned with time division multiple access and time division multiplexing in telecommunication networks. The invention is more particularly concerned with multirate frame structures, also known as multispeed frame structures, for tree-structured star telecommunication networks, especially fiber optic networks.\n2--Description of the Prior Art\nAs described in U.S. Pat. No. 5,363,370 by the inventor, a multirate frame is defined as a frame obtained by time division multiplexing of or time division multiple access to data having line bit rates, i.e., bit rates in one transmission medium, that are different; accordingly, a time interval, such as a frame time cell, having a constant duration can be allocated to a number of bits different from that allocated to another time interval according to the bit rates of the sources from which the bits originate. A multirate frame of this kind differs from a conventional frame produced by time division multiplexing of data at different bit rates, by increasing these different bit rates into a same line bit rate.\nIn prior art shared optic fiber networks there are two main modes of transmission between two equipments, namely half-duplex mode, also called as a two-way alternate transmission, and full-duplex mode, also called as a two-way simultaneous transmission, which results from wavelength division multiplexing.\nIn half-duplex mode the transmission direction alternates. While one equipment is sending data the other equipment is receiving the data, and vice versa. In each equipment there are a send phase and a receive phase that do not overlap. For a predetermined data bit rate D between the two equipments the line bit rate is equal to at least 2 D, given the alternating time sharing of the transmission medium between the two equipments. In practise a propagation time and an interference suppressing guard time increase the line bit rate to a value substantially equal to 2.5 D.\nIn wavelength division multiplexing full-duplex transmission, each of two equipments connected to respective ends of an optical fiber sends to the other equipment using a respective wavelength in the optical fiber.\nEach of these two transmission modes has its own advantages and disadvantages.\nHalf-duplex mode is particularly well suited to low bit rate data transmission. The doubling of the line bit rate does not impose any bit rate limitations on either the equipments or the transmission medium, and this technique can be implemented using reversible send/receive components which are not costly. In particular, a single diode can be used to send and receive by reversing the bias thereof. A diode of this kind has a range of operating frequencies enabling it to be used to send and receive at respective different frequencies.\nWavelength division multiplexing full-duplex mode is the only efficient transmission technique for high bit rate data. It uses costly optical duplexers to simultaneously enable reception at a first wavelength and transmission at a second wavelength."}
{"text": "1. Field of the Invention\nThe present invention relates to adaptive optics, and more specifically, it relates to a reference-free compensated imaging system for recovering a high resolution image of a target.\n2. Description of Related Art\nThe basic elements of a typical (prior art) adaptive optics system 100 are shown in FIG. 1. The goal of such systems is to provide real-time compensation for propagation errors, as encountered by optical beams, as they travel through dynamic distorting paths, including turbulent atmospheres, optical pointing errors, imperfect optical elements, multi-mode optical fibers, etc.\nBy compensating for optical wavefront distortions, one can enhance the performance of a variety of optical systems. Examples include optical communication systems, remote sensors, precision laser-beam delivery systems for industrial and medical purposes, and compensated imaging systems such as in medical applications (ophthalmological imaging and precision surgical procedures through the eye) and microscopy. In the latter example, this implies that one can view complex objects over a distorted path with the same image quality as if the path were distortion-free. In this case, the performance of the imaging system can approach that of its theoretical diffraction limit, within the so-called isoplanatic cone.\nIn what follows, we first discuss a generic adaptive optical system capable of correcting for path distortions encountered by a so-called reference beam. The reference beam is typically an image-free optical source, whose function is to sample the path distortions and, thus, provide this wavefront information as input to the adaptive optical system. This discussion is followed by a description of a specific adaptive optical configuration typical of prior-art, including an example of a wavefront-error sensing device. This, in turn, is followed by a discussion of an optical compensated imaging system typical of the art. An understanding of these prior-art systems will provide perspective with regard to the exemplary embodiments of this invention that follow.\nAs we discuss below, compensation of wavefront phase errors enables a system to provide diffraction-limited imaging and viewing of an extended object. In general, one first samples and compensates for propagation-path errors using, a diffraction-limited reference beam. Upon compensation of wavefront errors encountered by the reference beam, the optical system can approach its theoretical diffraction-limited imaging capability of image-bearing beams that lie within the so-called isoplanatic patch, which is well known in the art.\nIt is to be appreciated that the compensated optical imaging system can be implemented to service a variety of imaging based applications beyond atmospheric viewing systems. Hence, when the basic imaging system is referred to as a telescope, it is to be understood that the present teachings and embodiments can also be applied, without loss of generality, to compensated microscopy systems, speckle imaging, ophthalmological systems, communications systems, and the distortion path is referred to as a dynamic atmosphere, ultrasound imaging systems and so on. Similarly, when the distortion path that imposed the wavefront distortions to be compensated is referred to as a dynamic atmosphere, it is to be understood that the teachings can also be applied, without loss of generality, to correct for propagation-path distortions such as those experienced by imperfect optical elements, and static and/or dynamic distortions due to propagation through ocular systems, skin tissue, clouds, turbid liquids, and so on. The scene-based (Shack-Hartman) wave-front sensor could also be used in a post-processing scheme such as deconvolution or to augment speckle imaging.\nTurning now to FIG. 1, the goal of the prior art system is to enable one to view an optical source 110 with diffraction-limited capability. In this case, the optical source is chosen to be of spatial extent less than, or equal to, the diffraction limit of the optical system. Therefore, this source is equivalent to a point object with zero image-bearing information, analogous to a single pixel in an image. Light that emerges from this object, which is referred heretofore as a “reference beam,” 120, propagates through space, and, in general, becomes aberrated, as depicted by wavefront 120, as a result of the intervening path distortions or spatial phase errors, labeled by φ. In essence, the reference beam 120 samples the propagation path distortions between it and the optical compensation system, 100, including distortions imposed by optical elements within the compensation system itself.\nAt the receiver end of the link, a fraction of reference beam 120 is collected by telescope 130, which represents the input optical imaging elements of the adaptive optical receiver system 100. The collected light forms an image at the camera, or detector array, 190. In the absence of path distortions, the image at the camera plane would be in the form of an Airy disc, since the reference beam 120 is a sub-diffraction-limited point-source. However, owing to optical propagation phase distortions, φ, encountered by the reference beam on its path toward the receiver 110, the wavefronts of this beam will be aberrated, resulting in a distorted image of an Airy disc pattern at camera 190. As is known in the art, the path distortions in this scenario can stem from atmospheric turbulence, pointing and tracking errors, imperfect optical elements, thermal and mechanical perturbations, among other effects. The goal, therefore, of the adaptive optical system 100 is to compensate for such path errors so that the image quality of the reference beam at detector 190 can approach the diffraction limit.\nReturning to FIG. 1, the reference beam exiting the telescope 130 will be aberrated by virtue of the deleterious path distortions, as represented by wavefront 140. In this example, the adaptive optical system consists of two optical correction elements. The first corrective element 150 is a so-called tip-tilt compensator, whose function is to compensate for overall beam pointing and tracking errors. The second corrective element 160 is a spatial phase modulator, whose function is to compensate for fine-scale optical phase errors, including focus errors and spatially complex wavefront errors. The latter can include static and/or dynamic errors resulting from atmospheric turbulence and surface and volume refractive-index irregularities of optical elements, as discussed above. Wavefront compensation element 160 can be in the form of arrays of continuous and/or discrete optical phase shifters, such as piezoelectric transducers, electro-optic elements, deformable membranes, MEMS mirrors, liquid crystal cells, photonic crystals, among other devices, as is known in the art.\nThe incident distorted beam 140, first encounters the tip-tilt optical component 150 followed by the spatial phase modulator 160. The beam subsequently strikes a beam splitter 165, with one output beam directed to an optical wavefront error sensor 170, and with the other output beam directed to the camera detector 190.\nThe telescope provides an image of the incident beam at the camera plane 190, and, furthermore, provides an image of the pupil plane at the surface of the wavefront corrective element 160. Hence, the wavefront at the incident aperture is replicated, and scaled, as needed, at the plane of 160. The number of phase-controllable elements across the aperture of 160 is determined, in part, by the so-called transverse coherence parameter, otherwise known as the Fried parameter, which is characteristic of the scale size of the turbulent atmosphere.\nThe spatial bandwidth of the phase modulator 160 is designed to accommodate the spatial bandwidth indicative of the wavefront distortions, 120, subject to Nyquist constraints, as is known in the art. The sampling of the wavefront sensor 170 is also designed to accommodate the wavefront distortions 120 subject to Nyquist comstraints. In the image compensation systems (to be discussed with respect to FIG. 2 below), the spatial bandwidth requirements for the corrective plate are the same, in terms of resolving the wavefront error distortions sampled by the reference beam. The imaging resolution, on the other hand, is dictated by the diffraction limit of the overall optical system. In most cases, the Fried parameter scale size of the turbulence is far greater than that of the pixel size required to faithfully image the object.\nEach of the compensation elements 150 and 160 is controlled and configured in real-time using various classes of optical detectors, algorithms and electronic networks, examples of which are feedback, feed-forward and multi-dither systems, as is known in the art. One example of an optical feedback control loop is depicted in FIG. 1. It consists of a wavefront error sensor 170, a processor module 177, and a pair of electronic drivers 180 and 185 that provide control signals to tip-tilt compensator 150 and the spatial phase modulator 160, respectively. Ideally, the driver 185 will generate a spatial phase map indicative of a wavefront-reversed replica, whose phase is given by −φ. The resultant beam will therefore possess a wavefront that is a combination of the incident phase distortion, φ, with the correction phasemap, −φ, resulting in a wavefront with a net phase given as φ+(−φ)=0, indicative of an aberration-free reference beam.\nThe optical feedback control system is designed to correct the error 140 such that the wavefront at 190 is unaberrated. This is done by driving the wavefront error seen by 170 to a minimum relative to a known reference signal that encompasses the non-common-path errors of the optical system. Upon convergence of the servo control configuration, the resultant reference beam that strikes the camera/detector 190 will be, ideally, free of wavefront errors. In this state, the overall optical receiver system 100 will provide an image of the reference beam source 110, to its diffraction limit. Given that this system functions in real-time, dynamic path distortions can be tracked and compensated, with a residual error determined by the servo-loop gain and its bandwidth.\nTurning now to FIG. 2A, a compensated image adaptive optical system 200 is shown, typical of the prior art. The goal of this system is to enable precision imaging of an extended object 205 in the presence of dynamic path distortions 220, with the resultant image viewed by camera 290. The basic adaptive optical aspect of the system functions in a manner similar to that of FIG. 1. However, in the system depicted in FIG. 2A, there are now two different input beams incident upon a telescope 230. One of the two input beams is designated as a reference beam 110, and provides the same function as that of beam 110 of FIG. 1. That is, it is in the form of a sub-diffraction-limited optical source that samples the path distortions 220. The other incident light is an image-bearing beam of object 205 whose spatial information is also distorted by the path distortions 220, and whose high-fidelity compensated image is sought.\nThe reference and image-bearing beams both traverse the same input optical components and propagation path, including the telescope 230, intermediate focal plane 235, a collimation component, represented by lens 245, tip-tilt compensator 150, spatial phase modulator 160, imaging optics 247. The reference beam 110 and the image-bearing beam 205 both impinge upon beam splitter 265. The beam splitter directs each respective input beam into a different direction. The incident reference beam 110 emerges from one port of the beam splitter as beam 266 and propagates along one direction; and, the incident image-bearing beam 205 emerges from the other port of the beam splitter as beam 267 and propagates along a second direction. The reference beam 266 is directed to the adaptive optical control loop, and the image-bearing beam 267 is directed to a camera/detector module 290. Beam splitter 265 partitions the reference and image beams using a variety of discrimination techniques including polarization, wavelength, spatial frequency, temporal gating, as is known in the art.\nIn the compensated imaging system 200, the reference beam 266 exiting beam splitter 265 is directed to an adaptive optical processor in a manner analogous to that described with respect to FIG. 1. However, as opposed to FIG. 1, in the compensated imaging system of FIG. 2A, light from the incident reference beam 110 does not strike the camera 290. The sole purpose of the reference beam in this case is to provide path-distortion information to the wavefront error sensor 270 in the servo-loop so that, upon correction of the distortions imposed posed onto the reference beam, the image-bearing beam can be viewed with little or no distortion. The feedback loop, operationally, is similar to that of FIG. 1, namely, the raw wavefront-error information is inputted into processor 175 (see 277 in FIG. 2A), which, in turn provides error correcting information to drivers 180 and 185, the outputs of which provide signals to the tip-tilt compensator and the spatial phase modulator, 150 and 160, respectively.\nThe reference beam 266 emerging from beam splitter 265 passes through an intermediate image plane 255, followed by lens 249, which transforms the beam to a pupil plane. The beam is then scaled by the telescope (lenses 247 and 249) to satisfy the spatial bandwidth constraints of the wavefront-error sensor (WFS) 270. In this system, the WFS is a so-called Shack-Hartmann class of configuration. As is known in the art the Shack-Hartmann WFS consists of a lenslet array 271 and a detector array 273, the latter positioned at the focal plane of the lenslets. This pair of elements provides a spatial mapping of the local tilt phase errors across the overall pupil-plane aperture, that characterize the path-distorted incident reference wavefront 110. As known in the art, the required number of lenslets is a function of the square of the ratio of the input aperture size to that of the coherence (Fried) parameter indicative of the incident wavefront distortions. Under these constraints, it is assumed that the incident wavefront can be described as a series of plane-wave segments, each with a different tilt, or phase slope, and all concatenated together. Hence, each plane-wave segment is considered as a diffraction-limited beamlet, each with a different tilt angle.\nFIGS. 2B and 2C, respectively, illustrate the basic prior-art principles of the Shack-Hartmann WFS, as applied to an aberrated wavefront 220, and a distortion-free wavefront 221. The WFS, identical in both FIGS. 2B and 2C, consists of a lenslet array 271 and a multi-pixel detector array 273, the latter positioned at the focal plane of the lenslets. FIG. 2B depicts the operation of the WFS assuming an input reference beam whose wavefront is aberrated. Each plane-wave segment of the input beam 222 is incident upon a different lenslet in the array 271. In the presence of no wavefront phase errors beyond the wavefront sensor's Nysquist limit, a nearly diffraction-limited sinc-squared pattern will appear at each respective focal plane. However, since each plane-wave segment is comprised of a tilted wavefront, the sinc-squared pattern at each respective focal plane at the detector array 273 win be spatially shifted, with the lateral shift increasing with the slope of the local tilt. In most systems, especially atmospheric compensation systems, the spots will not form diffraction-limited pattern due to phase errors beyond. Nyquist which distort each individual spot. A “beam's eye view” at the detector surface 273, in the presence of the aberrated bean is shown in 274. Note that the array of focused spots is does not precisely overlap the grid-pattern. This is indicative of a typical aberrated beam, whose local tilts are randomly distributed. Therefore, each spot at the plane 274 has a correspondingly different offset in the (x,y) plane relative to the grid pattern. As is known in the art, the camera or ccd array 273 will require a sufficient number and density of resolvable detector pixels to measure the offset in spot position to ascertain the local tilt error with sufficient precision.\nFIG. 2C depicts the operation of the WFS assuming an input reference beam whose wavefront aberrations have been corrected. In the ideal case, the input beam 221 is a perfect plane wave, with a corresponding equi-phase surface across the entire input aperture to the WFS. As in FIG. 2B, each resolvable plane-wave segment of the input beam 223 is incident upon a different lenslet in the array 271. As before, an Airy disc pattern will appear at each respective focal plane along the detector surface 273. However, since each plane-wave segment has the same tilt (ideally, zero degrees with respect to the optical axis), each respective Airy pattern at the focal plane at the detector array 273 will be centered on its respective grid location. The “beam's eve view” at the detector surface 273, in the presence of the compensated reference beam, is shown in 274. Note that the array of focused spots precisely overlaps the grid-pattern. This is indicative of an ideal plane wave, whose local tilts are identical, and a wavefront sensor with no internal aberrations. In actual implementations, there will be internal wavefront sensor error that must be measured and removed to produce reference spot/sub-image default positions. Therefore, each spot at the plane 274 has a zero spatial offset in the (x,y) plane relative to the grid pattern. It is the goal of the servo-loop adaptive optical system to drive an aberrated beam (comprised of is finite number of tilted plane-wave segments) to a converged wavefront whose differential tilts approach zero.\nIt is important to emphasize that the WFS detects only the reference beam, which, by definition, does not contain image information, other than the spatial information resulting from the intervening propagation-path distortions. Hence, based on the prior art, in order to realize an image-compensation adaptive optics system, a reference beam must be present in addition to the image-bearing beam. However, in many applications, a diffraction-limited reference beam will not always be present or practical, even in cooperative scenarios (whereby, knowledge of the existence of a reference beam or of an observer is not a drawback). And, in certain cases, a reference beam optical source may be undesirable for strategic considerations, since detection of a reference optical source by a third party can reveal the presence and/or location of a covert observer. For these and other considerations, it is desirable to realize a compensated imaging system without the need for a cooperative reference beam"}
{"text": "This type of sensor may be remotely interrogated by connecting the input of the transducer to a radio frequency (RF) antenna. When the antenna receives an electromagnetic signal, the latter gives rise to waves on the surface of the substrate which are themselves reconverted into electromagnetic energy in the antenna. Thus, the device, consisting of a set of resonators connected to an antenna, has a response to the resonant frequencies of the resonators that it is possible to measure remotely. It is thus possible to produce remotely interrogable sensors. This possibility is a major advantage of surface acoustic waves and may be used notably in the context of pneumatic pressure sensors. This is because it is advantageous in this type of application to be able to place the sensor in the pneumatics whereas the interrogation electronics are stowed onboard the vehicle.\nAccording to the prior art, remote interrogation systems use interrogation signals in the form of pulses (typically with periods of about 25 μs) which are transmitted via an emitting antenna in the direction of a receiving antenna connected to the surface-wave sensor (referred to below, in the description, as a SAW sensor).\nA preferred frequency band for this type of system is the ISM (industrial, scientific and medical) band having a central frequency of 433.9 megahertz, the associated band width being 1.7 megahertz (MHz).\nGenerally, a remotely interrogable SAW sensor and its interrogation system may comprise, as illustrated in FIG. 1, in the simplified case of a single transducer:                an interrogation system 2; and        at least one resonator 1 comprising;                    an antenna 100; and            a comb transducer having interdigitated electrodes 11 and an SAW resonant cavity 13 characterized by its central frequency F and its quality factor Q (corresponding to the ratio between the central frequency and the width of the pass band). The cavity 13 comprises two series of reflectors that are uniformly spaced apart by a distance d. The transducer is connected to the antenna 100.                        \nThe interrogator 2 sends a long radio-frequency pulse so as to charge the resonator 1. After the emission has been stopped, the resonator discharges at its resonant eigenfrequency with a time constant τ equal to Q/πF. This discharge of the resonator forms the return echo detected by the receiver of the interrogator. Spectral analysis then allows the resonator frequency to be calculated and identified. This analysis may be carried out by algorithms based on Fourier transforms, for example an FFT (fast Fourier transform). This type of spectral-analysis treatment is particularly complex.\nA method for remotely interrogating passive SAW-type sensors has been proposed, in patent application WO 2008/015129, based on a frequency-modulation method. More precisely, this patent application discloses a method of measuring the resonant frequency of a resonator comprising the following steps:                emission, in succession, of radio-frequency signals of known carrier and modulation frequency, including the resonant frequency;        reception, via a receiving system, of response waves from the sensor; and        spectral analysis of the response waves from the sensor.        \nThe RF emission is frequency modulated with a modulation ωm and an amplitude modulation typically of about a kHz.\nThe response signal of the sensor is amplitude modulated with a modulation frequency ωm.\nThe measurement principle described above is based on the conversion of an (emitted) frequency modulation into an amplitude modulation via the transfer function of the resonator, as illustrated in FIG. 2. An emission modulated at an angular frequency ωm (corresponding frequency fm) results in the power detector receiving a signal modulated at the frequency ωm but more importantly, in a possible phase inversion of the modulated signal depending on whether it is below or above the resonant frequency.\nThe sinusoid, injected in the signal, at the frequency ωm, representing the modulating signal, is either directly converted into an amplitude modulation by the transfer function of the resonator (positive slope beneath the resonant frequency, amplitude modulation in phase with the frequency modulation), or inverted (negative slope above the resonant frequency, amplitude modulation in antiphase with the frequency modulation). The intermediate point corresponds to a null contribution to the received signal at the modulation frequency fm.\nAround this frequency position, the amplitude of the component of the signal at the frequency fm varies linearly.\nThe applicant observed that when implementing an algorithm providing the function described above, namely transformation of the contribution to the modulation frequency into a signed datum around the null contribution to the modulation frequency fm, and when detecting the amplitude, although a sinusoid is emitted to generate the frequency modulation, it is possible to use only two components of the received signal, respectively at the maximum and minimum modulation frequencies, this observation allowing interrogation times to be very substantially reduced."}
{"text": "1. Field of the Invention\nThe present invention relates in general to air conditioners, and relates in particular to an air conditioning system having a continuous air processing capability by alternately treating the process air through at least two desiccant members.\n2. Description of the Related Art\nFIG. 6 shows a prior art example of desiccant assisted air conditioning system same as the system disclosed in a U.S. Pat. No. 4,430,864. The system comprises: a process air passage A; a regeneration air passage B; two desiccant beds 103A, 103B; and a heat pump device 200 for desiccant regeneration and cooling of process air. The heat pump device 200 utilizes heat exchangers, embedded in the two desiccant beds 103A and 103B, as high and low temperature heat sources. In each of the thermal medium passages, there are opposingly disposed expansion valves 240A, 240B and one-way valves 241A, 241B, which are arranged parallel to the expansion valves 240A, 240B respectively, and the direction of compression of the compressor 230 can be switched by a four-way valve 250.\nIn the technology described above, cooling and dehumidifying processes can be explained with reference to a psychrometric chart shown in FIG. 7. The process air (state K) is withdrawn by a blower 102 through a passage 110, raised in pressure, and is forwarded to the one desiccant bed 103A through the passage 111 and the four-way valve 105 and passage 112A, where the moisture in the process air is adsorbed, to lower its humidity ratio and raise its temperature by the effect of the heat of adsorption. Because the desiccant bed 103A is cooled by the heat pump 200 through the heat exchanger 220, the adsorption heat is absorbed and the temperature of the process air does not rise too much, and after saturating (state L), the process air is dehumidified along iso-relative humidity line. The process air which has been dehumidified and maintained at the temperature (state N) is supplied to the conditioning space through the passage 113A, the four-way valve 106, passage 114. An enthalpy difference .DELTA.Q is thus produced between the return air from the conditioning space (state K) and the cooled process air (state N), to provide cooling of the conditioning space.\nThe regeneration process of the desiccant is performed as follows. Regeneration air (state Q) is withdrawn into the blower 140 through the passage 120, raised in pressure, and is forwarded to the other desiccant bed 103B through the passages 121, 122, the four-way valve 106, and the passage 113B. The desiccant bed 103B is heated by the heat pump 200 by way of the heat exchanger 210, so its temperature is raised, and the relative humidity is lowered (state R). The regeneration air which now has a lowered relative humidity passes through the desiccant bed 103B to remove the moisture from the desiccant material (state T). The regeneration air which has passed through the desiccant bed 103B passes through the passage 112B, four-way valve 105 and the passage 124 and is discharged to an outside environment.\nAfter the air conditioning process has been carried out for sometime and the moisture content in the desiccant becomes higher than a certain value, the four-way valve is operated to be switched, so that the air passages for the desiccants and cooling/heating of the heat pumps are interchanged. Thus, the operation is carried on so that the regenerated desiccant is used to continue air conditioning operation while the other desiccant is being regenerated. Therefore, it can be seen that the processes of adsorption and regeneration are conducted in a batch type system.\nIn the technology described above, heat exchange of the low temperature heat source of the heat pump and the desiccant for adsorption are embedded into a unit, and heat exchange of the high temperature heat source of the heat pump and the desiccant on the regeneration side are embedded into a unit. So, the cooling effect .DELTA.Q is provided by a direct thermal load on the heat pump (refrigeration device), which means that it is not possible to generate more cooling than that allowed by the capacity of the heat pump acting as a refrigeration device. Therefore, this configuration does not provide any advantages worthy of making the apparatus complex. In addition, there has been required two four-way valves, one for reversing the operation cycle of the heat pump and the other for interchanging the passages of the process/regeneration air, which further makes the configuration of the apparatus complex."}
{"text": "This invention relates to a brake especially adapted for use with all-terrain vehicles or over-the-road towed vehicles, such as boat trailers or the like. More specifically, this invention relates to such a brake wherein the components thereof are sealed from the environment.\nOver-the-road trailers and all-terrain vehicles are normally provided with conventional single disc caliper brakes with, at times, being provided with drum brakes. As such, the components of these brakes are exposed to the atmosphere which readily results in the corrosion of many of their component parts. Such is particularly true in the case of boat trailers where the brakes are frequently exposed to water, even salt water. Such corrosion not only affects brake wear, efficiency, and consistency, but also will ultimately lead to the need to replace many critical, and expensive, brake components. Moreover, heat dissipation problems with these prior art brakes can affect brake performance and potentially lead to premature failure.\nWhen a brake is sealed to the atmosphere, usually the components thereof are immersed in oil for cooling purposes. However, the oil must be sealed from the brake fluid used to operate the brakes of over-the-road trailers or all-terrain vehicles, which presents a problem of seal compatibility; that is, the seal material used for petroleum-based oil is incompatible with the brake fluid or vice-versa. The solution of the prior art to this problem is to provide two seals with an air gap between them which allows each fluid to interact only with the seal with which it is compatible. Such undesirably adds to the size and the cost of the brake.\nIn addition to the foregoing problems, the prior art brakes are difficult and time consuming to install and/or replace. With respect to trailers, first, a spindle must be welded onto the axle of the trailer and bearings must be pressed onto the spindle. Then a hub and disc, or hub and drum, assembly, with bearing races formed therein, must be attached to the spindle. Then a bracket must be mounted on the frame of the trailer and a caliper attached to the bracket adjacent to the disc. The wheel is then mounted to the hub. With respect to all-terrain vehicles, the caliper and wheel assembly must be mounted to the strut on the front of the vehicle. First, bearings must be installed into bores cut into the strut itself. Then, a shaft is installed into the bearings. The shaft is splined on the outside to accept the gear tooth configuration of a hub which mounts both the wheel and the rotor for the caliper. The caliper is mounted to a machined face of the strut by a separate bracket which allows the caliper to float on the disc.\nIn view of these problems, it is evident that the need exists for a modular brake which is easy to install and which will isolate all of the components of the brake from the deleterious effects of its environment.\nIt is thus an object of the present invention to provide a brake for a trailer or all-terrain vehicle which is modular in nature and which totally houses the brake components to seal them from the environment.\nIt is another object of the present invention to provide a brake, as above, which is easy to install.\nIt is a further object of the present invention to provide a brake, as above, in which the brake components are protected from corrosion.\nIt is an additional object of the present invention to provide a brake, as above, which provides consistent braking forces for the vehicle.\nIt is yet another object of the present invention to provide a brake, as above, which provides superior heat dissipation and otherwise prevents premature failure.\nIt is a still further object of the present invention to provide a brake, as above, in which the components thereof are immersed in oil and a single seal can be utilized between the activating brake fluid and the oil.\nIt is still a further object of the present invention to provide a brake, as above, in which the components will not disadvantageously wear, thereby providing a brake which will not normally need replacement or service.\nThese and other objects of the present invention, as well as the advantages thereof over existing prior art vehicle brakes, which will become apparent from the description to follow, are accomplished by the improvements hereinafter described and claimed.\nIn general, a brake for a wheel of a vehicle made in accordance with the present invention includes a housing, and an axle assembly adapted to be attached to the wheel of the vehicle for rotation therewith. A plurality of stationary discs are positioned in the housing. A plurality of rotatable discs are carried by the axle assembly and are positioned to be engaged by the stationary discs. At least one bracket is carried by the housing. Each stationary disc has at least one ear which is received in the bracket. A piston is moveable to cause the stationary discs to engage the rotatable discs to apply a braking torque to the wheel.\nIn accordance with another aspect of the present invention, the axle assembly includes drive surfaces which engage a shaft sleeve having a splined outer surface. The plurality of rotatable discs are carried by the splined outer surfaces.\nIn yet another aspect of the present invention, a first fluid is provided in the housing. The piston includes a ring and a rib and is moveable by a second fluid to cause the ring to engage a stationary disc to cause the stationary discs to engage the rotatable discs to apply a braking torque to the wheel. A seal is provided adjacent to the rib to prevent the first fluid from communicating with the second fluid.\nIn a further aspect of the present invention, the housing has open ends, and a plate is attached to one end of the housing and is adapted to be attached to the vehicle. A seal is provided between the plate and the housing. The axle assembly is received in the other end of the housing and a seal is provided between the housing and the axle assembly.\nA preferred exemplary vehicle brake incorporating the concepts of the present invention is shown by way of example in the accompanying drawings without attempting to show all the various forms and modifications in which the invention might be embodied, the invention being measured by the appended claims and not by the details of the specification."}
{"text": "1. Field of the Invention\nThe present invention relates generally to methods and systems for creating and managing product and project portfolios. The present invention also relates to methods and systems for creating, evaluating, ranking viewing potential returns of proposed products and projects. The present invention further relates to methods and systems for enabling decision makers to view the contributions proposed products would make over their lifetime and the cash flow that would result from investment and returns. The present invention also relates to methods and systems to enable decision makers to view the consumption of key resources. Also, the present invention relates to methods and systems for developing project or product scenarios based on selected assumptions.\n2. Description of the Related Art\nProduct development today poses unprecedented challenges. To stand out and succeed in competitive and crowded marketplaces, products must be continuously innovative and fresh. As short-lived market opportunities present themselves, product development efforts that once required years now must be executed in a matter of months.\nCompanies must continuously test new ideas. Scarce investment funds must be allocated to the most promising ideas and opportunities. Projects must be tracked to ensure that the funds committed to those projects are not better employed elsewhere. Product development departments are typically many times overcommitted. Product managers and the executive management team often have no way of determination whether the project under development is well aligned with corporate objectives. Indeed, there currently exists no unified and easy way of determining where the company's finite resources are committed. Moreover, today's executive teams have no way of effectively balancing the risks already extant in the company's investment portfolio with the risks inherent in proposed or existing product development efforts or proposals. Management has no way of seeing an objective view of aggregate resource load, nor do they have the systems to make staffing and resource decisions other than through gut instinct or the opinion of the person requesting such resources. This provides no basis for reallocation decisions or staffing decisions that must therefore flow up the organization as a “take it or leave it” lump. What are needed, therefore, are methods that would provide an analyst or executive the opportunity to balance the profit impact of making changes to the resource pool with postponing or cancelling profitable projects. What are also needed are methods of enabling such analysts to see which resources that may be undercommitted as well as those that may overcommitted, to yield opportunities for profitable rescheduling and to manage the risks in those bottlenecks.\nFrom the foregoing, it may be appreciated that a new strategic planning tool is needed that would assist and provide a unified framework for decision makers in the strategic planning process. Such a tool should lower the cost of the planning cycle, result in a faster time to market and lower the risk of new investment proposals.\nTo lower the cost of the planning cycle, such a tool should automate the collection of assumptions on which an investment proposal is based, including the collection of financial, resource and marketing information. Other desirable features of such a tool include the ability to rank proposals across many investment criteria and the ability to select those proposals that best for the company's existing portfolio. Also desirable is the ability to create scenarios that include different mixes of investment projects, products and funding and the ability to see the resultant returns of such mixes, preferably summarized across financial criteria and non-financial criteria, such as manpower, risk and marketing criteria, for example. To achieve a faster time to market, product managers and development executives should have at their disposal a tool that enables them to judiciously allocate the limited resources available to them to the most profitable projects, and to enable them to allocate the right resources to the right opportunities. To lower the risks of new projects, the inherent risks of proposed investment and new projects should be calculated as a measure of how accurate the assumptions need to be and till produce a positive return. Preferably, such a tool should present investment proposals in a uniform and balanced manner, both in terms of short term and long-term investments, thereby lowering the risk of being overtaken by competitors."}
{"text": "When partitioning a memory, it is conventional to partition the memory in contiguous sectors such that the logical mapping and the physical mapping are identical. Data is then stored in the sectors of the memory."}
{"text": "With the completion of the sequencing of the human genome, one of the next grand challenges of molecular biology will be to understand how the many protein targets encoded by DNA interact with other proteins, small molecule pharmaceutical candidates, and a large host of enzymes and inhibitors. See e.g., Pandey &\nMann, “Proteomics to study genes and genomes,” Nature, 405, p. 837-846, 2000; Leigh Anderson et al., “Proteomics: applications in basic and applied biology,” Current Opinion in Biotechnology, 11, p. 408-412, 2000; Patterson, “Proteomics: the industrialization of protein chemistry,” Current Opinion in Biotechnology, 11, p. 413-418, 2000; MacBeath & Schreiber, “Printing Proteins as Microarrays for High-Throughput Function Determination,” Science, 289, p. 1760-1763, 2000; De Wildt et al., “Antibody arrays for high-throughput screening of antibody-antigen interactions,” Nature Biotechnology, 18, p. 989-994, 2000. To this end, tools that have the ability to simultaneously quantify many different biomolecular interactions with high sensitivity will find application in pharmaceutical discovery, proteomics, and diagnostics. Further, for these tools to find widespread use, they must be simple to use, inexpensive to own and operate, and applicable to a wide range of analytes that can include, for example, polynucleotides, peptides, small proteins, antibodies, and even entire cells.\nBiosensors have been developed to detect a variety of biomolecular complexes including oligonucleotides, antibody-antigen interactions, hormone-receptor interactions, and enzyme-substrate interactions. In general, biosensors consist of two components: a highly specific recognition element and a transducer that converts the molecular recognition event into a quantifiable signal. Signal transduction has been accomplished by many methods, including fluorescence, interferometry (Jenison et al., “Interference-based detection of nucleic acid targets on optically coated silicon,” Nature Biotechnology, 19, p. 62-65; Lin et al., “A porous silicon-based optical interferometric biosensor,” Science, 278, p. 840-843, 1997), and gravimetry (A. Cunningham, Bioanalytical Sensors, John Wiley & Sons (1998)).\nOf the optically-based transduction methods, direct methods that do not require labeling of analytes with fluorescent compounds are of interest due to the relative assay simplicity and ability to study the interaction of small molecules and proteins that are not readily labeled. Direct optical methods include surface plasmon resonance (SPR) (Jordan & Corn, “Surface Plasmon Resonance Imaging Measurements of Electrostatic Biopolymer Adsorption onto Chemically Modified Gold Surfaces,” Anal. Chem., 69:1449-1456 (1997), (grating couplers (Morhard et al., “Immobilization of antibodies in micropatterns for cell detection by optical diffraction,” Sensors and Actuators B, 70, p. 232-242, 2000), ellipsometry (Jin et al., “A biosensor concept based on imaging ellipsometry for visualization of biomolecular interactions,” Analytical Biochemistry, 232, p. 69-72, 1995), evanascent wave devices (Huber et al., “Direct optical immunosensing (sensitivity and selectivity),” Sensors and Actuators B, 6, p. 122-126, 1992), and reflectometry (Brecht & Gauglitz, “Optical probes and transducers,” Biosensors and Bioelectronics, 10, p. 923-936, 1995). Theoretically predicted detection limits of these detection methods have been determined and experimentally confirmed to be feasible down to diagnostically relevant concentration ranges. However, to date, these methods have yet to yield commercially available high-throughput instruments that can perform high sensitivity assays without any type of label in a format that is readily compatible with the microtiter plate-based or microarray-based infrastructure that is most often used for high-throughput biomolecular interaction analysis. Therefore, there is a need in the art for compositions and methods that can achieve these goals."}
{"text": "The present invention relates to a method of producing a zipper tape agglutinated in use to a bag for packaging as well as to a zipper tape produced by the production method, and more specifically to a zipper tape having a mating section comprising a male strip and a female strip and a method of producing the zipper tape.\nBags for packaging which have a zipper tape closure are very popular. These bags can easily and repeatedly be opened and closed. It is generally expected in the market that a bag with a zipper tape will be used more and more in the future.\nThe conventional type of zipper tape is generally based on a polyethylene system or a polypropylene system, and a male strip and a female strip forming a pair extending in a row and used to form a mating section. The strips are integrally formed on a tape, and the tape is cut in use along a center line thereof to provide a slender section with a male strip and that with a female section. When this zipper tape is agglutinated in use to a bag 1 for packaging as shown in FIG. 1, the cross section taken along the line Axe2x80x94A in the figure is as shown in FIG. 2(A), and a tape section 4 with a male strip 3a and a female strip 3b formed thereon is based on a polyethylene system when a resin used for an innermost layer 5 of a bag 1 for packaging is a polyethylene-based one, while the tape section 4 is based on a polypropylene system when the resin for the innermost layer 5 is a polypropylene-based one.\nHowever, the conventional type of zipper tape has the defects as described below.\na) At first, when the zipper tape 2 is produced by means of miscellaneous shape extrusion molding, the male strip 3a and the female strip 3b are respectively formed on the tape section 4 in an integrated state therewith as shown in FIG. 2(A), and in this production method, it is impossible to mold a plurality of male strips and a plurality of female strips simultaneously in batch, and the two types of strips are formed only as a pair by means of extrusion molding, so that the productivity is low.\nb) When the zipper tape 2 is agglutinated to the bag 1 for packaging, when the innermost layer 5 of the bag 1 and the tape section 4 of the zipper tape 2 are jointed to each other by loading a pressure and heating the jointed section from the outside as shown in FIG. 2(A), as there is the risk that the tape section 4 itself may be fused there is no way but to make the thickness of the tape section 4 larger.\nc) Because of the requirement described above, when a pressure and heat are applied from the outside for making side seals 1a, 1b of the bag 1 for packaging as shown in FIG. 2(B), the mating section 3 and the tape section 4 form together a fused resin block 3xe2x80x2, 4xe2x80x2, and the section expands, so that the pressure and heat required for sealing must be strong, which causes other sealed sections to be heated excessively. On the contrary, when the pressure and heat used for sealing are weak, the resin block 3xe2x80x2, 4xe2x80x2 is not fully fused such that pin holes are easily generated therein.\nd) In the zipper tape 2 comprising the mating section 3 made from a polyethylene-based resin and the tape section 4, as the innermost layer 5 of the bag 1 for packaging is thermally fused only the polyethylene-based bag, so that the zipper tape 2 can not be used for a bag made from a polypropylene-based resin. On the other hand, the zipper tape 2 made from a polypropylene-based resin can not be used for the bag 1 made from a polyethylene-based resin, and therefore the zipper tape 2 must be selected according to a material used for forming the innermost layer 5 of the bag 1.\nImproved types of zipper tapes have been disclosed for the purpose to solve the problems as described above. As the most representative technology, there is the technology in which a heat-insulating film 6 having a melting point higher than that of the tape section 4 is laminated on the tape section 4 of the zipper tape 2 and further a film 7 for thermal fusion having a melting point lower than that of the heat-insulating film 6 is laminated on the heat-insulating film 6 (Japanese Patent No. 1908600). Namely, in the zipper tape 2, which is applied to the bag 1 for packaging shown in FIG. 1, the cross section Axe2x80x94A is as shown in is enlarged form in FIG. 3(A), while the cross section Axe2x80x2xe2x80x94Axe2x80x2 is as shown in its enlarged form in FIG. 3(B).\nWith this improved type of zipper tape 2, some of the problems in the conventional technology have been solved, but the following ones are still left unsolved.\ne) When the zipper tape 2 is produced by means of the miscellaneous shape extrusion molding, the male strip 3a and the female strip 3b are formed integrally on the tape section 4 respectively, so that there is no way but to form the two types of strips in a pair and in a row, so that the productivity is not improved.\nf) When the innermost layer 5 of the bag 1 and the thermally-fused film 7 of the zipper tape 2 are pressed to each other and heated for sealing, the risk that the tape sections 4 facing to each other are thermally fused to each other is evaded because of the presence of the heat-insulating film 6 as shown in FIG. 3(A), and this improvement advantageously makes it possible to reduce the total thickness of the tape section 4, heat-insulating film 6, and thermally fused film 7 down to about 100xcexc.\ng) When the side seals 1a, 1b are made in the bag 1, the thickness of the fused resin block 3xe2x80x2, 4xe2x80x2, where the mating section 3 and the tape section 4 are integrated to each other, is reduced in proportion to the reduced thickness of the tape section 4 as shown in FIG. 3(B), but the reduction in thickness is still insufficient.\nh) In a case of the zipper tape 2 made from a polyethylene-based resin, also the thermally fused film 7 and the innermost layer 5 of the bag 1 must be made from the polyethylene-based resin, and further in a case of the zipper tape 2 made from a polypropylene-based resin, also the thermally fused film 7 and the innermost layer 5 of the bag 1 must be made from the polypropylene-based resin, and this technological limitation still remains unsolved.\ni) In the conventional type of zipper for a bag, two sheets of base tapes to be sealed on an inner surface of the bag, and the male strip and the female strip respectively formed on this base tape are integrally formed by an extrusion molding machine.\nBecause of this requirement, when it is desired to give a specific color to the zipper tape for the purpose to indicate that the zipper tape is attached to inside of an upper open section of the bag or to make the zipper tape colorful, it is necessary to add a coloring material to a resin used as a material for the tape and to extrude the colored resin from the extrusion molding machine.\nj) For the reason described above, it is possible to give a specific color to the zipper tape as a whole, but it is impossible to give a color only to the male and female strips, nor to give different colors to the mating section and base tape respectively, and therefore when designing a bag with a zipper, options in designing the zipper section are disadvantageously limited as compared to those allowable in other sections.\nk) Further, once the zipper tape has been set in an upper open section of a bag, only a surface of the base tape and the mating section can directly be checked visually, and it is impossible, because of the presence of the mating section, to print on the zipper tape an alerting message such as, for instance, that xe2x80x9cIf the zipper tape color has not changed to blue, please check and confirm again that the zipper tape has completely be closedxe2x80x9d for alerting a user to check whether the mating section has been completely closed. It is desired to print this type of alerting message there.\nThe present invention was made in the light of the circumstances as described above, and it is a first object of the present invention to provide a zipper tape with the large width and having a structure allowing formation of a plurality of pairs of a male strip and a female strip each pair extending in a row on the zipper tape and a method of producing the zipper tape.\nIt is a second object of the present invention to provide a zipper tape making it possible to reduce a fused resin block, which is generated at the mating section when side sealing for the bag is made, as much as possible to prevent generation of pin holes therein and also to ensure the provision of a good appearance of the bag and also to provide the method of producing the same.\nFurther it is a third object of the present invention to provide a zipper tape compatible with an innermost layer of a bag regardless of whether the innermost layer is made from a polypropylene-based resin or a polyethylene-based resin and also capable of eliminating the necessity of selecting a zipper tape made from the same material as that used for an innermost layer of the bag and also to provide the method of producing the zipper tape.\nIt is a fourth object of the present invention to provide a zipper tape for a bag allowing printing of any letter and figure or the like with any color on a surface of the base tape.\nIt is a fifth object of the present invention to make it possible for a user of a bag to visually determine that a zipper tape has been fitted to inside of an upper open section of the bag by printing thereon a remark indicating presence of the zipper tape or giving any specific color to the zipper tape.\nIt is a sixth object of the present invention to improve commercial value of a product packaged in a bag for the purpose to contribute to increase of sales in the market by making it possible to make a design of a surface of the bag more colorful especially in the section where the zipper tape is fitted to the bag."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for adjusting the brightness level of an image and an image processing apparatus and an image processing program for adjusting the brightness level of the image based on the method.\n2. Description of the Related Art\nGenerally, the range of a brightness level which can be reproduced on photographic paper (photosensitive material) or on a monitor is narrower than that of a brightness level which can be obtained by an imaging element of a digital camera. Therefore, a function for adjusting the brightness level of the image is provided in an apparatus which reproduces a photographic image obtained by the digital camera on the photographic paper or the monitor.\nFor example, an apparatus and a method for outputting a print of an image after compressing the dynamic range of the brightness level of the image is disclosed in U.S. Pat. No. 6,122,076. In this apparatus, processing is performed so that a sharp image is output even if the dynamic range of the image is compressed. When the brightness level of the image is adjusted so that the difference in lightness/darkness is reduce, it is important that sufficient contrast is maintained so that a fine-detail structure in the image does not become unsharp.\nPhotography by an electronic endoscope is performed by emitting light which is guided by a light guide to a body cavity, which receives no natural light. Since the body cavity has a narrow complicated shape, light unevenly illuminates the body cavity. Therefore, the difference in lightness/darkness of an image obtained by endoscopy is much larger than that of a photograph of a person, which is taken against light. Further, the color of the body cavity is substantially the same through the entire region. Therefore, it is difficult to recognize a fine-detail structure in an image obtained by endoscopy compared with a photograph obtained by a digital camera or the like. Further, in endoscopy, it is necessary to reproduce a fine-detail structure in a region which is used for diagnosis. However, the region which is used for diagnosis is not always a dark region.\nThe method disclosed in U.S. Pat. No. 6,122,076 may be also applied to image processing of an electronic endoscope apparatus. However, the photograph obtained by the digital camera and the image obtained by endoscopy are different from each other, as described above. Therefore, a more sophisticated method for adjusting a brightness level, which is appropriate for image processing of the electronic endoscope apparatus, is desired."}
{"text": "The present invention relates to an image input unit and an image input method for inputting a color image to be printed out.\nHitherto, in a printer, complementary color ink has been used to finally output image data representing the values of CMY (Cyan, Magenta, Yellow) or CMYK (Cyan, Magenta, Yellow, black) and in a monitor, image data representing the value of RGB (Red, Green Blue) is finally output. Generally, in image processing using a computer, color matching is performed with the display colors of a monitor as the reference and therefore in an image input unit such as a digital still camera, image data is represented in the area of RGB color space with the RGB values taking 0 or more. For example, in not only a digital still camera using primary color CCD, but also a digital still camera using complementary color CCD, when CMYG (Cyan, Magenta, Yellow, Green) image data provided from an input image is converted into RGB image data, processing of adjusting the components of a conversion matrix or replacing negative values with 0 is performed and only RGB values of 0 or more are output.\nFor example, in sRGB standard, color coordinates of RGB are determined and RGB color space wherein the RGB values are 0 or more is defined. However, the sRGB color space is narrower than the color space that can be recognized by the visual sense of a human being and thus only a color space in a narrower range than the visual sense of a human being can be represented in an image input unit for outputting an image represented in the sRGB color space.\nOn the other hand, in a printer, a color space wherein the CMY or CMYK values are 0 or more can be represented by printing. A printer for outputting an image represented in the color space wherein the CMY or CMYK values are 0 or more can use pigments for color development on the outside of the sRGB color space as CMY ink, thereby representing a wider color space than the sRGB-defined color space by printing. Likewise, the printer can represent a wider color space by printing than a color monitor manufactured with sRGB display as the reference.\nFor such a printer to print image data recorded in an image input unit such as a digital still camera for recording an image represented in a color space wherein RGB colors represented in sRGB are 0 or more, the image is represented only in a narrower range than the color representation area of the printer, and the characteristics of the printer are not used effectively; this is a problem.\nIn recent years, image data format standards allowing negative values in the colors of sRGB64, RIMM-RGB, etc., have been developed. However, if an attempt is made to allow negative values and improve the image quality, the information amount per pixel is increased three bits and thus the image data formats result in an increase in the data capacity."}
{"text": "When processing molten metal such as aluminum, dissolved gases including hydrogen are removed through degassing techniques. Metals including magnesium are removed from the molten metal through demagging techniques. Gases commonly used for degassing are nitrogen and argon while gas commonly used for demagging is chlorine.\nWhen processing molten aluminum, a bath of the molten aluminum is contained within a vessel such as a furnace. A layer of dross is formed at the surface of the molten metal bath composed of various chemical compounds such as aluminum oxide, magnesium compounds, flux and refractory particles. Some particles formed of magnesium compounds are too small to float to the surface of the bath. Chlorine gas released into the molten aluminum bath bonds with such magnesium compounds forming magnesium chloride, which can be removed from the surface of the bath. Solid fluxes of various compositions can be added to the molten metal to remove magnesium or other impurities. For example, U.S. Pat. No. 3,650,730 discloses adding solid flux containing chlorine salt instead of using chlorine gas. Adding solid flux to molten metal is dangerous because workers are near the molten metal when the flux is added. Chlorine gas is extremely toxic and chlorine gas that does not react with magnesium in the bath may enter the surrounding area creating a hazardous workplace.\nThe processing of molten metal commonly employs pumps that include various components depending on the application, including circulation, transfer and gas purification pumps. A gas purification pump disclosed in U.S. Pat. No. 5,993,728 to Vild, is used for injecting chlorine gas into molten metal to react with magnesium such as from aluminum can scrap. The pump includes a submerged base having an inlet opening leading to an interior impeller chamber. A discharge passageway leads from the impeller chamber to an exterior of the pump. An impeller is rotated in the impeller chamber, which draws molten metal through the inlet into the impeller chamber and out the discharge passageway. The chlorine gas is injected into the discharge passageway.\nIn conventional practice as shown by U.S. Pat. No. 4,052,199 to Vild, solid flux can be manually added to molten aluminum in another chamber downstream of the pump to remove magnesium prevalent in aluminum can scrap.\nU.S. Pat. No. 6,589,313 discloses a hollow shaft on the end of which is an impeller. The shaft and impeller are rotated and positioned at an angle relative to the bath by a complex apparatus. Solid flux and gas is added to the rotating shaft and dispersed in the molten metal.\nA device for degassing molten metal without pumping includes a motor powered vertical shaft and rotor on the end positioned in a well or location of a furnace. A source of pressurized gas is connected to the shaft. The shaft and rotor include a passageway for the gas. The spinning of the rotor disperses the gas in the molten metal.\nAn improved device for injecting gas into a pump for pumping molten metal is disclosed in published U.S. patent application Pub. No. 2006/0180962.\nAn improvement disclosed in U.S. patent application Ser. No. 11/691,664, filed Mar. 27, 2007, entitled “FLUX INJECTION WITH PUMP FOR PUMPING MOLTEN METAL,” features a machine for feeding solid flux into a gas vortex for entraining the solid with the gas and directing it into molten metal discharged from a pump."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a dielectric film in a capacitor. In particular, the present invention relates to a method of manufacturing a dielectric film having a crystallization preventing layer between two metal-oxide layers by way of an atomic deposition layer process. The inventive method provides a dielectric film exhibiting excellent electrical properties, while achieving cost-effective throughput.\n2. Description of the Related Art\nIn general, capacitors may have a structure of Metal-Insulator-Metal (MIM) to improve their electrical properties. Among such capacitors, the structure of Tin/HfO2/Tin, as compared, for example, to the structure of Ru/TaO2/Ru, may exhibit low leakage current and high capacitance. As such, capacitors having the structure of Tin/HfO2/Tin may have an advantage of being applied to semiconductor memory devices with a design rule of 100 nm or less. However, the Tin/HfO2/Tin capacitors may have a relatively high leakage current during their back-end processing as a result of the HfO2 layer crystallization.\nTypically, crystallization of the HfO2 layer during the back-end processing of the capacitors may be prevented by interposing an Al2O3 layer between two HfO2 layers as a crystallization preventing layer to form a dielectric film of the structure HfO2/Al2O3/HfO2. The relatively high crystallization temperature of the Al2O3 layer may increase the effective crystallization temperature of the overall dielectric film, thereby reducing leakage current during back-end processing even at high processing temperatures.\nA dielectric film composed of HfO2/Al2O3/HfO2 may be formed by an Atomic Layer Deposition (ALD) process to provide a uniform layer deposition on a surface of a lower electrode of a capacitor. Such uniform step coverage and thickness of a dielectric film on the electrode may require increased deposition time of precursor materials; a factor which may result in overall increased processing time of the dielectric film and reduced throughput of the process.\nAn attempt has been made to increase the throughput of the ALD process, while maintaining good product quality, by employing a batch-type manufacturing process, i.e., depositing films on more than one wafer at a time. However, when the deposition time of the required precursor materials is increased, the resulting overall throughput in such a process is almost indistinguishable over an ALD processing of a single wafer. On the other hand, when the deposition time of the precursor materials is decreased, the precursor materials are not properly applied to the lower electrode of the capacitor, thereby reducing the step coverage of the film and providing inferior wafer thickness and film uniformity, consequently producing an electrically defective capacitor.\nAccordingly, there exists a need for a method of manufacturing a dielectric film by an ALD process that produces high quality dielectric films, having excellent step coverage and electrical properties, while providing cost-effective throughput."}
{"text": "The present invention relates to the processing of a signal in order to determine jitter properties of the signal.\nJitter is a major cause for bit errors in digital communication systems. Jitter describes timing aberrations of a signal from a predefined signal course, e.g. delays of a signal caused by a lossy transmission channel. Bit errors occur if a receiver of a digital communication system detects a logical zero in a received bit stream instead of a logical one or vice versa. A measure for bit errors in digital communication systems is the bit error ratio—BER.\nParticularly, the BER depends on the quality of a received signal containing digital data. Receivers of digital communication systems are adapted to sample the received signal with a predetermined sampling rate. However, jitter of the received signal may cause a bad sampling of the digital data transmitted with the received signal causing bit errors.\nThe theoretical model of total jitter—TJ—distinguishes between random jitter—RJ—and deterministic jitter—DJ. RJ is caused by thermal and noise effects. Thus, it is unbounded and may be described by a Gaussian probability density function. In contrast to RJ, DJ is bounded, i.e., it has definite amplitude limits and is by nature a peak-to-peak value. DJ is further distinguished in periodic jitter—PJ—, bounded uncorrelated jitter—BUJ—, and data dependent jitter—DDJ.\nDDJ is a function of bit patterns and further distinguished in Duty Cycle Distortion—DCD—and Inter-Symbol Interference—ISI. ISI is usually caused by long and short bit cycles and results from bandwidth limitations or from loss within transmission lines while DCD is caused by voltage offsets between differential inputs and differences between transition times in a system.\nUntil now Jitter decomposition is known using a Real Time or Sampling Scope, as described in the publication “Jitter Analysis: The dual-Dirac Model, RJ/DJ, and Q-Scale”-White Paper, 21 Jun. 2005.\nState of the Art for Bit Error Ratio Testers—BERT—is separating RJ and DJ using a BERT scan methodology. Such BERT might also be used to measure TJ, using a so-called Optimized BERT Scan Method."}
{"text": "This invention relates to utensil supporting grates of cooktops and, more particularly, to a cooktop/grate combination which provides greater stability of the grates.\nFor many years grates were mounted to cooktops by inserting depending fingers or knobs of the grate through mating openings in the cooktop. This had a number of drawbacks. For example, it required that holes extend through the cooktop, providing sites for possible rusting of the body of the cooktop.\nMore recently the burner openings have been surrounded by relatively large recessed areas and the grates have been mounted within these recessed areas. In one embodiment the grate includes an annular base with a vertical outer surface and a bottom edge. The recessed area of the cooktop includes spaced apart bosses with horizontal upper surfaces that support the bottom edge of the grate base. The recessed area also includes dams with vertical inner surfaces that contact the outer surface of the grate base to restrain the grate should a lateral force tend to slide the grate out of the recessed area.\nAt times such dams do not exert sufficient restraining force on the grates, which then can slide out of the recessed area. A significant reason for such sliding is that the dams do not have sufficient vertical overlap with the grate base. This results from the restraints on the vertical dimensions of various elements.\nGrates must support heavy utensils when at high temperatures. Thus they are made of relatively large amounts of expensive materials. Thus it is desirable to minimize the size of the grate. At the same time, for optimum operation of the burner flame, the grate must support the utensil at a predetermined distance above the burner. For this reason, the grate does not rest on the bottom of the recessed area. Rather bosses are formed part way up the peripheral wall of the recess and the grate is supported on the bosses. This enables the base of the grate to have a smaller height and use less material.\nIn order to restrain the grates from sliding out of the recessed areas substantially vertical opposing surfaces are provided. Essentially the entire periphery of the grate base can be used as one vertical surface, while the other must be formed along the periphery of the recessed area. The recessed areas usually are large in cross-section and deep. Attempts to draw such a recessed area with a generally vertical wall into the steel or other metal forming the cooktop leads to excessive failures of the metal. Therefore, individual small dams, with substantially vertical inner surfaces, have been formed in spaced apart relationship about the peripheral portion of the recessed area. These dams and the corresponding portions of the grate base form the opposing vertical surfaces. In some instances the area so provided is not sufficient.\nOne possible solution would be to remove or shorten the bosses. However, since the top of the grate must be at least a predetermined distance above the burner, this would require greater height of the grate base and would greatly increase the cost of each grate.\nAn object of the present invention is to provide an improved cooktop and grate assembly.\nAnother object is to provide such an assembly in which the grates are resistant to sliding out of the mating recess in the cooktop.\nIt is yet another object of this invention to provide such an improved cooktop and grate assembly that does not involve any increase in material."}
{"text": "It is often necessary in semiconductor processing to fill high aspect ratio gaps with insulating material. This is the case for shallow trench isolation (STI), inter-metal dielectric (IMD) layers, inter-layer dielectric (ILD) layers, pre-metal dielectric (PMD) layers, passivation layers, etc. As device geometries shrink and thermal budgets are reduced, void-free filling of narrow width, high aspect ratio (AR) features (e.g., AR>6:1) becomes increasingly difficult due to limitations of existing deposition processes."}
{"text": "1. Field of the Disclosure\nEmbodiments disclosed herein relate generally to threaded connections. More particularly, embodiments disclosed herein relate to two-step wedge thread connections and related methods of makeup.\n2. Background Art\nCasing joints, liners, drill pipe, and drill collars (collectively referred to as “tubulars”) are often used in drilling, completing, and producing a well. Casing joints, for example, may be emplaced in a wellbore to stabilize a formation, to protect a formation against elevated wellbore pressures (e.g., wellbore pressures that exceed a formation pressure), and the like. Casing joints may be coupled in an end-to-end manner by threaded connections, welded connections, and other connections known in the art. The connections may be designed so as to form a seal between an interior of the coupled casing joints and an annular space formed between exterior walls of the casing joints and walls of the wellbore. The seal may be, for example, an elastomeric seal (e.g., an o-ring seal), a metal-to-metal seal formed proximate the connection, or similar seals known in the art. In some connections, seals are formed between the internal and external threads. Connections with this characteristic are said to have a “thread seal.” As used herein, a “thread seal” means that a seal is formed between at least a portion of the internal thread on the box member and the external thread on the pin member.\nIt will be understood that certain terms are used herein as they would be conventionally understood where tubular joints are being connected in a vertical position along a central axis of the tubular members such as when making up a pipe string for lowering into a well bore. Thus, the term “load flank” designates the side wall surface of a thread that faces away from the outer end of the respective pin or box member on which the thread is formed and supports the weight (i.e., tensile load) of the lower tubular member hanging in the well bore. The term “stab flank” designates the side wall surface of the thread that faces toward the outer end of the respective pin or box member and supports forces compressing the joints toward each other such as the weight of the upper tubular member during the initial makeup of the joint or such as a force applied to push a lower tubular member against the bottom of a bore hole (i.e., compressive force). The term “face” of the box is the end of the box member facing outward from the box threads and the term “nose” of the pin is the end of the pin member facing outward from the threads of the connection. Upon makeup of a connection the nose of the pin is stabbed into and past the face of the box.\nOne type of thread commonly used to form a thread seal is a wedge thread. In FIGS. 1A and 1B, a connection 100 having a wedge thread is shown. “Wedge threads” are characterized by threads that increase in width (i.e., axial distance between load flanks 125 and 126 and stab flanks 132 and 131) in opposite directions on the pin member 101 and box member 102. Wedge threads are extensively disclosed in U.S. Pat. No. RE 30,647 issued to Blose, U.S. Pat. No. RE 34,467 issued to Reeves, U.S. Pat. No. 4,703,954 issued to Ortloff, and U.S. Pat. No. 5,454,605 issued to Mott, all assigned to the assignee of the present invention and incorporated herein by reference.\nOn the pin member 101, the pin thread crest 122 is narrow towards the distal end of the pin member 101 while the box thread crest 191 is wide. Moving along the axis 105 (from right to left), the pin thread crest 122 widens while the box thread crest 191 narrows. The rate at which the threads change in width along the connection is defined by a variable known as the “wedge ratio.” As used herein, “wedge ratio,” although technically not a ratio, refers to the difference between the stab flank lead and the load flank lead, which causes the width of the threads to vary along the connection. Furthermore, as used herein, a thread “lead” refers to the differential distance between a component of a thread on consecutive threads. As such, the “stab lead” is the distance between stab flanks of consecutive thread pitches along the axial length of the connection. In FIGS. 1A and 1B, the thread surfaces are tapered, meaning that the pin thread 106 increases in diameter from beginning to end while the box thread 107 decreases in diameter in a complimentary manner Having a thread taper improves the ability to stab the pin member 101 into the box member 102 and distributes stress in the connection.\nFor wedge threads, a thread seal is accomplished by the contact pressure caused by interference over at least a portion of the connection between the pin load flank 126 and the box load flank 125 and between the pin stab flank 132 and the box stab flank 131, which occurs when the connection is made-up. Close proximity or interference between the roots 192 and 121 and crests 122 and 191 completes the thread seal when it occurs over at least a portion of where the flank interference occurs. Higher pressure may be contained with increased interference between the roots and crests (“root/crest interference”) on the pin member 101 and the box member 102 and by increasing flank interference. This particular connection also includes a metal-to-metal seal that is accomplished by contact between corresponding sealing surfaces 103 and 104 located on the pin member 101 and box member 102, respectively.\nA property of wedge threads, which typically do not have a positive stop torque shoulder on the connection, is that the make-up is “indeterminate,” and, as a result, the relative position of the pin member and box member varies more for a given torque range to be applied than connections having a positive stop torque shoulder. As used herein, “make-up” refers to threading a pin member and a box member together. A final make-up refers to threading the pin member and the box member together up to a desired amount of torque, or based on a relative position (axial or circumferential) of the pin member with the box member.\nFor wedge threads that are designed to have both flank interference and root/crest interference at a selected make-up, both the flank interference and root/crest interference increase as the connection is made-up (i.e. increase in torque increases flank interference and root/crest interference). For wedge threads that are designed to have root/crest clearance, the clearance decreases as the connection is made-up. Regardless of the design of the wedge thread, corresponding flanks and corresponding roots and crests come closer to each other (i.e. clearance decreases or interference decreases) during make-up.\nIndeterminate make-up allows for the flank interference and root/crest interference to be increased by increasing the torque on the connection. Thus, a wedge thread may be able to thread seal higher pressures of gas and/or liquid by designing the connection to have more flank interference and/or root/crest interference or by increasing the torque on the connection; however, this also increases stress on the connection during make-up, which could lead to failure during use.\nPrior to make-up a flowing joint compound commonly referred to as “pipe dope” is typically applied to surfaces of a threaded connection to improve the thread seals and provide lubrication during make-up of the connection. For example, the base (e.g., a grease) of the pipe dope may assist a wedge-threaded connection in achieving a thread seal between load and stab flanks thereof, e.g., as disclosed in U.S. Pat. No. RE 34,467 issued to Reeves. Further, pipe dope may contain metallic particle additives, such as copper to protect the threads of the pin and box members from friction galling during make-up and break-out.\nWhen a wedge thread connection is made-up, excess pipe dope may become trapped (rather than being squeezed out) between engaging pin and box threads, which may either cause false elevated torque readings (leading to insufficient make-up or “stand-off”) or, in certain circumstances, damage the connection. Pipe stand-off due to inadequate evacuation of the pipe dope is detrimental to the structural integrity of wedge thread connections. As the pressure build-up may bleed off during use, the connection is at risk of accidentally backing off during use. Therefore, stand-off in wedge thread connections is of particular concern as it may lead to loss of seal integrity or even mechanical separation of the two connected members.\nFIG. 2 shows a prior art two-step connection 150. The threads that form the connection are separated on two different “steps,” a large step indicated by the bracket 31 and a small step indicated by the bracket 32. The portion between the large step 31 and the small step 32 is commonly referred to as a mid-step 901. In some connections, the mid-step 901 may be used as a metal-to-metal seal. The pin thread crest 222 on the small step 32 of the pin member 101, at its full design height, does not interfere with the box thread crest 221 on the large step 31 of the box member 102 when the pin member 101 is stabbed into the box member 102. The diameter of the small step 32 of the pin member 101 is smaller than the smallest crest-to-crest thread diameter on the large step 31 of the box member 102. The pin thread 106 on the small step 32 can be stabbed past the box thread 107 on the large step 31. The threads on both the small step 32 and the large step 31, which have substantially the same nominal lead, engage with each revolution to make-up the connection. Thus, the number of revolutions during which the threads slide or rub against each other is reduced for the same number of engaged threads. A two-step connection allows for each of the steps to have threads with different characteristics as long there is little or no variance in the nominal lead of the threads on the steps.\nA two-step wedge thread connection is disclosed in U.S. Pat. No. 6,206,436 issued to Mallis and assigned to the assignee of the present invention. That patent is incorporated herein by reference. Mallis discloses a two-step wedge thread connection having different wedge ratios, one of which is considered to be an “aggressive” wedge ratio and the other a “conservative” wedge ratio. “Aggressive” refers to the larger wedge ratio, and “conservative” refers to the smaller wedge ratio. Everything else the same, the greater the wedge ratio, the more determinate the make-up. Too large of a wedge ratio may have an inadequate wedging effect, which can allow the connection to back-off during use. Smaller wedge ratios are better able to resist backing-off of the connection. Too small of a wedge ratio may have such an indeterminate make-up that galling may occur over the lengthened make-up distance. Mallis discloses that one of the steps can have a wedge ratio that is optimized for a more determinate make-up (aggressive), while the other step can have a wedge ratio that is optimized for preventing back-off of the connection (conservative).\nFIGS. 3A and 3B show cross-section views of a conventional two-step wedge thread connection 200 prior to a final makeup. The connection 200 includes a pin member 201 having pin wedge threads 106 thereon and a box member 202 having corresponding box wedge threads 107 thereon. Further, the connection 200 has first step 31 and second step 32, with a mid-step region 901 located therebetween. As shown, an axial separation of the two wedge thread steps 31, 32 of the pin member 201, indicated by distance ‘B,’ is substantially equal to an axial separation of the two wedge thread steps 31, 32 of the box member 202, indicated by distance ‘A.’ Thus, the pin wedge threads 106 on the first step 31 of the pin member 201 may be characterized as “in-phase” with the box wedge threads 107 on the first step 31 of the box member 202. Likewise, the pin wedge threads 106 on the second step 32 of the pin member 201 may be characterized as “in-phase” with the box wedge threads 107 on the second step 32 of the box member 202.\nThe corresponding pin and box threads on the two steps 31, 32 are in-phase such that during makeup, gaps 137 between approaching load flanks 125, 126 are equal to gaps 138 between approaching stab flanks 131, 132 on the first step 31. Similarly, gaps 137 between approaching load flanks 127, 128 are equal to gaps 138 between approaching stab flanks 133, 134 on the second step 32. Thus, corresponding load flanks 125, 126 and stab flanks 131, 132 on first step and corresponding load flanks 127, 128 and stab flanks 133, 134 on second step 32 will contact at substantially the same time (i.e., at final makeup). FIGS. 3C and 3D illustrate the conventional two-step wedge thread connection 200 at a final makeup. Because the corresponding load flanks 125, 126 and stab flanks 131, 132 on first step 31, and corresponding load flanks 127, 128 and stab flanks 133, 134 contact at substantially the same time, the interference generated between the surfaces is substantially equal.\nA threaded connection having improved make-up and break-out torque characteristics would be appreciated by those skilled in the art."}
{"text": "Flexible pipe is useful in a myriad of environments, including in the oil and gas industry. Flexible pipe may be durable and operational in harsh operating conditions and can accommodate high pressures and temperatures. Flexible pipe may be bundled and arranged into one or more coils to facilitate transporting and using the pipe.\nCoils of pipe may be positioned in an “eye to the side” or “eye to the sky” orientation. When the flexible pipe is coiled and is disposed with its interior channel facing upwards, such that the coil is in a horizontal orientation, then the coils of pipe are referred to as being in an “eye to the sky” orientation. If, instead, the flexible pipe is coiled and disposed such that the interior channel is not facing upwards, such that the coil is in an upright or vertical orientation, then the coils of pipe are referred to as being in an “eye to the side” orientation.\nThe flexible pipe may be transported as coils to various sites for deployment (also referred to as uncoiling or unspooling). Different types of devices and vehicles are currently used for loading and transporting coils of pipe, but usually extra equipment and human manual labor is also involved in the process of loading or unloading such coils for transportation and/or deployment. Such coils of pipe are often quite large and heavy. Accordingly, there exists a need for an improved method and apparatus for loading and unloading coils of pipe."}
{"text": "1. Technical Field of Industrial Application\nThe present invention relates to the technology of applying a high-viscosity bonding agent having a small thickness to a member.\n2. Related Art\nIn the case of an ink-jet recording head, for example, which has a nozzle plate with a plurality of nozzle holes bored therein, a pressure generator, an ink supply port, a spacer for partitioning a common ink chamber and a second plate, which are fixedly formed in layers so that ink drops are made to jet out of the nozzle holes under the pressure generated in a pressure generating chamber by means of piezoelectric vibrators, piezoelectric vibrators whose tips are approximately tens to hundreds of .mu.m in size and which are arranged at intervals of tens of .mu.m must be fastened to a vibrating plate by means of a bonding agent.\nIn order to fasten such microscopic members disposed at small intervals to the vibrating plate with a bonding agent, the bonding agent, which is as viscous as possible in an unhardened state, needs to be applied uniformly without overcoating at the time of pressure-bonding.\nFor the reason stated above, a bonding agent is spread on a board with a squeegee or the like as thin as possible and subsequently the bonding agent thus spread on the board is transferred and spread onto an adherent as disclosed in Japanese Unexamined Patent Publications Nos. 91274/1982 and 150974/1992. However, the minimum possible thickness of such a bonding agent that can be spread by a squeegee or the like is in the order of tens of .mu.m, which is far greater than that of an ideal bonding agent layer for use in the manufacture of ink-jet recording heads. Such a large thickness causes a decrease in the production yield which arises from overcoating and the like at the time of bonding."}
{"text": "Heart disease continues to take a heavy toll of human life. It is responsible for nearly a million deaths a year in the U.S. alone. The toll has been reduced to some extent by improved and new medical procedures, new drugs, pacemakers and most recently by the employment of ventricular assist devices. However, in many cases, replacement of the diseased heart with a healthy heart is the only manner of treating and curing the patient. One of the major problems with this method of treatment is that there are too few donors. For example, in 1978 only 31 people received transplanted hearts, whereas a panel of the National Heart, Lung and Blood Institute has estimated that from 17,000 to 50,000 people per year could use a transplanted heart.\nA possible solution to this problem is a totally artificial heart. The development of such a device began in 1957. Since 1975 a number of artificial hearts have been developed and tested in animals. Unfortunately, most of the artificial hearts developed prior to 1970 were not successful for various reasons. In 1970 a series of developments made the prospect for the device look better. A major advance was the design of a heart whose pumping element was a diaphram. This principle was introduced by Clifford S. Kwan-Gett. Basically, the artificial heart developed by Kwan-Gett comprises a chamber containing a deformable diaphram. The heart was driven by an external pump which periodically pumps fluid into the heart chamber to deform the diaphram and thereby expel blood from the artificial heart chamber. Notwithstanding, the improved diaphram system, Kwan-Gett's artificial heart, along with other hearts of similar design and construction, did not prove to be as successful as had been hoped due to shape and size problems.\nIn 1972, the diaphram principle was used by Robert Jarvik to design a more anatomically acceptable artificial heart. In animal testing, the Jarvik artificial heart has been very successful. The most recent model, the Jarvik-7 heart which was developed in 1979 has been remarkably successful. The Jarvik-7 heart has been implanted in calves which subsequently have survived for considerable periods of time. Despite the success of the Jarvik heart, the system utilized to power the heart remains a major problem which has stymied the development of artificial hearts. The pneumatically powered artificial hearts, including the Jarvik-7 heart, that have proven successful in animals are not portable. The animal is confined to a cage, tethered to a large drive system and exercised only on a tread mill. Such conditions would be unacceptable for human beings. Even if compressed air devices were made portable, the large pneumatic tubes that enter the chest would be uncomfortable and would carry a high risk of infection at the point of entry.\nVentricular assist devices, which assist the natural heart by supplementing the function of the left ventricle have spurred the development of an improved power system. While not as large as artificial hearts, ventricular assist devices work on the same principle as the artificial heart. Thus, they too require a pneumatic or hydraulic power source. An example of one such ventricular assist device is disclosed in Runge U.S. Pat. No. 4,176,411. The device disclosed by Runge comprises a flexible dacron conduit disposed between the left atrium of the heart and the descending thoracic aorta and surrounded by an artificial muscle sheath. In operation, the artificial muscle sheath is electrically stimulated to contract against the flexible dacron conduit to force blood to flow from the left atrium in to the aorta, thus aiding corporeal circulation.\nThe most recent development in an improved power system is an electrohydraulic energy converter that is presently being design by Jarvik in collaboration with Milton Isaacson of New-Tech, Inc. This electrohydraulic energy converter has only one moving part. The impeller of an axial flow pump is attached to the rotor supported by a single hydrodynamic bearing. Reversing the rotation pump reverses the direction of the hydraulic flow. The hydraulic fluid (silicon oil of low viscosity) actuates the diaphram of a blood pump much as compressed air does. In a left ventricle cardiac assist device the axial flow pump moves the hydraulic fluid from a reservoir sack into the blood pump and back. In a total artificial heart the hydraulic fluid is pumped back and forth between the right and left ventricles. The energy converter is so small that it can be implanted without impinging on vital structures. The converter, however, requires an external battery and an electronics package which is connected to the heart by a small cable that passes through the chest. The batteries weigh from 2 to 5 lbs. and are worn in a vest or on a belt. It is necessary to replace the rechargable batteries once or twice a day. Although this device contains many of the desirable features of a portable artificial heart for human use, it nevertheless has many of the problems of the totally external pneumatic pumps. The electrohydraulic energy converter is not a completely internal power system, but rather is dependent upon an external power source. Also, like the external pneumatic pumps, this new device has wires that enter the chest which would result in a high risk of infection at the point of entry. Moreover, this new device is dependent on a battery system which must be continually recharged and which must be worn on a belt or vest.\nThe present invention solves the problems of the prior art pumping devices since it is directed to a totally internal system for powering artificial hearts, ventricular assist devices and the like.\nThe present invention was first disclosed to the U.S. Patent and Trademark Office in Disclosure Document 97,636 dated Feb. 10, 1981."}
{"text": "A solar cell array has groups of solar cells (hereinafter called “cell string”) in which a plurality of solar cells are connected in series, bypass diodes (hereinafter called “CST bypass diode”) that are connected to each of the groups in parallel, and reverse-current preventive diodes that are connected to one or more cell strings in series.\nFIG. 1 is an illustration of the circuit of a solar cell array. A solar cell array 91 contains a cell string 92a in which a plurality of cells (95a1, 95a2, and 95an) are connected in series and a cell string 92b in which a plurality of cells (95b1, 95b2, and 95bn) are connected in series in the same way. CST bypass diodes 93a and 93b are connected to the cell strings 92a and 93b in parallel, respectively. In addition, CST bypass diodes 93c and 93d are connected to other cell strings 92c and 92d in parallel, respectively. A reverse-current preventive diode 94p is connected to the cell strings 92a and 92b in series. A reverse-current preventive diode 94q is connected to the cell strings 92c and 92d in series in the same way. Bypass diodes for a single cell (96a1, 96a2, 96an, 96b1, 96b2, and 96bn) are placed to protect each cell from reverse current. Said bypass diodes for a single cell may be omitted.\nIn conventional approaches, there have been proposed solar cell modules in which a CST bypass diode, together with a plurality of solar cells, is integrally sealed between the front surface member and the rear surface member of the modules. The CST bypass diode for the proposed solar cell modules comprise a bear chip diode, an attaching section having conductive members arranged on each side of the bear chip diode (such conductive members constitute a mounting surface), and a base unit made from materials with excellent heat dissipativity, which base unit is arranged together with the attaching section outside the bear chip diode (refer to the patent document 1, for example).\nFurthermore, a diode in which the upper surface and the under surface of the diode with lead terminals connected are arranged almost on the same plane has been proposed as a bypass diode for a single solar cell (refer to the patent document 2, for example).                Patent document 1: Japanese Patent Laid-Open No. 2000-164910 bulletin        Patent document 2: Japanese Patent Laid-Open No. 2000-277785 bulletin        "}
{"text": "This invention relates to a process and an apparatus for manufacturing a cast product and, more particularly, to a cast forging process and apparatus capable of manufacturing a machine component to which a high-strength is required, such as a scroll of a scroll compressor and a VTR drum without any fear of defect such as shrinkage cavities.\nFIGS. 7a and 7b are schematic sectional views illustrating a conventional indirect insertion-type squeeze casting apparatus disclosed in \"Illustrated Dictionary of Casting, Second Edition\" P. 121, edited by Japanese Association of Casting, FIG. 7a illustrating a state before casting and FIG. 7b illustrating a state immediately after casting. In these figures, reference numeral 1 is a top die, 2 is a bottom die, 2a is a molten metal reservoir disposed in the bottom die 2, 2b and 2c are a plurality of gates and cavities (a single cavity may be used), respectively, disposed between the top and the bottom dies 1 and 2, 3 is an extrusion plunger, 4 is a die-clamping ram attached to the top die 1, 5 is a knock-out punch, 6 is a molten metal and 7 is a solidified casing.\nThe operation will now be described. The top die 1 is firmly urged against and fastened to the bottom die 2 by the die-clamping ram 4. The molten metal 6 is supplied to the molten metal reservoir 2a formed in the bottom die 2 immediately before the casting. In this state, the injection plunger 3 is moved downward to pressurize the molten metal 6. The injection plunger 3 continues to pressurize the molten metal 6 at a high pressure equal to or more than 500 atmospheric pressure until the molten metal 6 solidifies and casting is completed. Then, the injection plunger 3 is moved rearward (upward) and the top die 1 is moved upward by the die-clamping ram 4 to open the die and the cast product 7 can now been taken out from the die by moving the knock-out punch 5 upward. The gate sections 2b are removed by a separate pressing step to obtain a cast product corresponding to the cavity portion 2c.\nFIGS. 8a and 8b illustrate a scroll member for use in a scroll compressor having a hollow boss portion, FIG. 8a being a front view and FIG. 8b being a sectional view taken along B--B line in FIG. 8a. In these figures, reference numeral 40 is a scroll member for a scroll compressor, which has a scroll-shaped spiral teeth portion 40a on one side of a base plate portion 40b and a boss portion 40c axis-symmetrically formed on the other side of the base plate portion 40b. This is a component which requires a high precision of the order of .mu.m.\nFIG. 9 illustrates a schematic sectional view of an indirect insertion cast forging apparatus for casting the scroll member for the scroll compressor illustrated in FIGS. 8a and 8b based on the disclosure of \"Illustrated Dictionary of Casting, Second Edition\" P. 121, edited by Japanese Association of Casting, the left half illustrating the state before casting and the right half illustrating after casting. In the figure, reference numeral 1c is a gate disposed in the top die 1 and 1d is a cavity.\nThe operation will now be described. The top die 1 is firmly urged against and fastened to the bottom die 2 by the die-clamping ram 4. The molten metal 6 is supplied to the molten metal reservoir 2a formed in the bottom die 2 immediately before the casting. In this state, the injection plunger 3 is moved upward to pressurize the molten metal 6. The molten metal 6 is introduced by the injection plunger 3 to the cavity 1d through the gate 1c. The injection plunger 3 continues to pressurize the molten metal 6 at a high pressure equal to or more than 500 atmospheric pressure until the molten metal 6 solidifies and casting is completed. Then, the die-clamping ram 4 is moved-upward to open the die 1, and the injection plunger 3 is moved upward and a connected cast product in the cavity 1d and the molten metal reservoir 2a can now been taken out from the die. The gate sections 1c are removed by a separate step of sawing to obtain a cast product 40 corresponding to the cavity portion 1d.\nAs is well known, the cast forging process is a process in which a molten metal is injected into a die cavity at a low speed capable of maintaining a laminar flow to prevent the generation of cast cavities due to the trapped gas and in which a high pressure about 500 atmospheric pressure or more is kept being applied to the molten metal until it is solidified to prevent the generation of shrinkage cavities due to the volume shrinkage upon the solidification of the molten metal, the process being known as a process for obtaining a defect-free, quality cast product. However, while it is necessary to keep applying a high pressure until each corner of the molten metal in the die solidify in order to prevent the generation of the shrinkage cavities, in the case of the indirect insertion cast forging process illustrated in FIGS. 7a and 7b, the molten metal 6 flows through the gate 2b to reach the cavity 2c and solidifies therein, so that it is necessary that the molten metal within the gate 2b does not solidify and stays in a molten state until the molten metal within the cavity 2c solidifies in order that the high pressure applied to the molten metal 6 within the molten metal reservoir 2a from the injection plunger 3 keeps being transmitted to the molten metal within the cavity 2c. Also, the molten metal within a recessed portion at the center of the cavity 2c in the illustrated example is required to solidify at a later time than the molten metal surrounding it. That is, in order to prevent the generation of the shrinkage cavity, it is necessary that a directional solidification in which the solidification is achieved progressively from the outer portion of the cavity 2c toward the molten metal reservoir 2a can be achieved. In order to achieve this, it is necessary that the gate 2b has large dimensions, which leads to problems of requiring a machine cutting of the gate 2b to obtain the cast product corresponding to the cavity 2c after the casting 7 is taken out from the die because the press cutting provides a very poor surface conditions and dimensional accuracy of the cut surface. The second problem is the limited configuration of the product, in which the dimension of the thin-wall portion for example is limited because the molten metal which becomes a product in the cavity 2c must be solidified from the outer portion toward the inner portion.\nFurther, in order to prevent generation of shrink cavities in the scroll member of the scroll compressor illustrated in FIGS. 8a and 8b, it is necessary that the directional solidification in which the solidification of the molten metal progresses from the outer portion of the cavity 1d toward the molten metal reservoir 2a can be achieved as in the apparatus illustrated in FIG. 9. In order to achieve this, firstly, it is necessary that the gate 1c has large dimensions, which leads to problems of requiring a machine cutting of the gate 1c to obtain the cast product corresponding to the cavity 1d after the casting is taken out from the die because the press cutting provides a very poor surface conditions and dimensional accuracy of the cut surface. Secondly, with this configuration of the product, the shrink cavities in the base plate portion near the gate cannot completely be eliminated even when the gate thickness is made similar to the thickness of the base plate portion."}
{"text": "Wet capacitors are increasingly being used in the design of circuits due to their volumetric efficiency, reliability, and process compatibility. Wet capacitors typically have a larger capacitance per unit volume than certain other types of capacitors, making them valuable in high-current, high power and low-frequency electrical circuits. One type of wet capacitor that has been developed is a wet electrolytic capacitor that includes a valve metal anode, a cathode, and a liquid electrolyte. The unit cell voltage in this type of capacitor is generally higher due to the formation of a dielectric metal oxide film over the anode surface. Wet electrolytic capacitors tend to offer a good combination of high capacitance with low leakage current. Another type of wet capacitor is a wet symmetric capacitor in which the anode and cathode are similar in terms of structure and composition. The unit cell voltage in this type of capacitor is generally low due to the inevitable decomposition of the electrolyte at high voltage. Whether electrolytic or symmetric, however, the cathodes of wet capacitors typically include a substrate and a coating that provides high capacitance through a faradic or non-faradic mechanism (e.g., double layer formation). Conventional coatings include activated carbon, metal oxides (e.g., ruthenium oxide), and the like. Unfortunately, however, the coatings can become easily detached under certain conditions, such as in the presence of aqueous electrolytes.\nAs such, a need remains for a high voltage wet electrolytic capacitor that possesses good mechanical robustness and electrical performance."}
{"text": "1. Field of the Invention\nThe present invention relates to pressure transducers and, more particularly, the present invention is directed to a disposable transducer apparatus for use in an electromanometry system for monitoring and recording hemodynamic pressures.\n2. The Prior Art\nWith the advent of high speed digital computers in the 1950's and 1960's, a number of medical teams throughout the United States began considering the possibility of computerized blood pressure monitoring. Today, computerized blood pressure monitoring systems are frequently used in intensive care facilities to provide around-the-clock computer monitoring of hemodynamic pressures in critically ill heart patients both during and following cardiovascular surgery and the like.\nAs the heart pumps, the resistance to fluid flow and elasticity of the blood vessels cause a buildup in the arteries of a continuous mean blood pressure having a superimposed fluctuation between a maximum and minimum pressure, known as systolic and diastolic pressures, respectively. It is the function of a computerized blood pressure monitoring system to continuously measure and electronically record measurements of the cyclical blood pressures which occur by virtue of the heart's pumping action.\nIn order to transform these measurements of cyclical blood pressures into recordable electronic waveforms, a catheter filled with a sterile saline solution is inserted into a patient's artery and is threaded through the artery until the tip of the catheter is located near the heart. In this way, the periodic pulsations of blood are transmitted as mechanical pulses of fluid through the saline-filled catheter. The catheter tubing is connected to a pressure transducer which converts the mechanical pressure pulses through the saline in the catheter into electronic waveforms that may be visually displayed on an oscilloscope screen or on a recording strip chart. These electronic waveforms are subsequently analyzed by a computer to determine a number of quantitative parameters, including heart rate, duration of systole, and systolic, diastolic and mean pressures. This information gives the surgeon and other health care personnel a quantitative look at the operation of a patient's heart, and is thus extremely valuable in helping to assess the clinical course of critically ill patients on a continuously updated basis.\nThe pressure transducers typically used in an electromanometry system are relatively expensive. Thus, in the past the practice has been to resterilize the transducer after using it on a patient so that it could be used anew on other patients. This practice tended to increase the risk of cross-contamination among patients, and also resulted in increased maintenance costs and inconvenience. Moreover, because such transducers are expensive, if one was accidentally damaged or dropped, the expense to the hospital became significant. A further disadvantage of the prior art transducers is their size and weight. Because of these size and weight considerations prior art pressure transducers were mounted on brackets attached to a standard I.V. pole. This separate mounting caused inconvenience and undue complication when setting up the electromanometry system or transporting the patient. The prior art type transducers also require frequent calibration, thus adding to the maintenance costs and inconvenience already present with these devices.\nIn an effort to reduce the expense and inconvenience associated with having to resterilize non-disposable transducers after each use, disposable pressure transducer domes have been developed which physically isolate the transducer from the saline solution carried by the catheter tubing. The disposable transducer domes typically consist of a thin plastic membrane which is used to physically separate the diaphragm of the transducer from the saline solution so that the diaphragm is not contaminated by the saline solution, which has been in contact with the patient's blood. Although such disposable transducer domes help to minimize the risk of contamination and also help to reduce the cost associated with resterilization of pressure transducers, the plastic membrane attenuates the mechanical pressures transmitted through the saline solution, with the result that the resulting electronic waveforms are often degraded in quality.\nFurther efforts to improve the state of the art concerning pressure transducers used in connection with blood pressure monitoring have included efforts to develop smaller, miniaturized versions of the pressure transducer. The smaller, miniaturized pressure transducers, although much reduced in size so that they can be directly supported on the arm of the patient or by the pressure tubing, nevertheless are expensive to construct and are still intended to be reused, thus requiring continued use of disposable transducer domes.\nThus, what is needed in the art is a very small, lightweight transducer which is sufficiently economical in its construction as to be disposable after each use, thus eliminating the need for resterilization techniques and also eliminating the need for using disposable transducer domes."}
{"text": "1. Field of the Invention\nThis invention relates to medical IV administration line connectors. More particularly, this invention relates to needleless injection ports for the safe infusion and/or aspiration of fluids in intravenous and blood administration therapy.\n2. State of the Art\nIntravenous therapy has a long history of use in supplying patients with pharmaceuticals, liquid nourishment, or blood products. Prior art FIGS. 1 through 4 show that the current or conventional way of introducing parenteral liquid solutions and/or blood and blood products into a patient is by the conventional gravity feeding system 10. The feeding system 10 includes a container 12 that is either a bottle or bag for the parenteral solution, a tube 14 extending from the bottle or bag and connected to a Y-injection site 16 (piggyback or secondary Y-injection site), and a tube 18 from the Y-injection site 16 to a needle or catheter 20 which is inserted into a vein in the arm 22 of the patient. The vein-access needle or catheter is taped to the patient with adhesive tape 24 so that the chance of a detachment and disconnect from the vein is minimized.\nSupplemental intravenous therapy from a piggyback or secondary bottle or bag 26 is introduced through the Y-injection site 16 into the primary intravenous administration set 10. As seen best in FIG. 3, the Y-injection site 16 which is integrated into the primary intravenous administration set 10 consists of two tubular conduits 16a, 16b which merge into a third tubular conduit 16c. The tubing 12 from the bottle or bag of parenteral solution of the primary intravenous administration set 10 is attached into the inlet port 16a of the Y-injection site. In similar fashion, the tube 18 is attached into the exit port 16c of the Y-injection site. A sealed entry port segment 17 of the Y-injection site 16 is provided by the extension conduit 16b which has a standard, self-sealing latex rubber septum 17a at its inlet port to seal this port from leakage. Consequently, it is difficult for pathogens to enter the Y-injection site 16 via the septum port 17 because of the seal 17a. This self-sealing septum 17a is of a conventional design and includes coaxial annular aprons which fit over the conduit wall and grip the external and internal wall surfaces to hold the septum securely to the conduit 16b. Typically, a plastic shrink-band (not shown) is shrunk on the outer wall of the septum 17a to securely connect it to the extension conduit 16b.\nThe supplemental intravenous solution is introduced into the primary intravenous administration set 10 through the Y-injection site 16 by way of a primed piggyback or secondary intravenous set 26. The piggyback or secondary intravenous set 26 has a hollow-bore needle 28 attached to its distal end, which in turn is inserted through the self-sealing septum 17a of the Y-injection site 16 and into the extension conduit 16b. This needle 28 is connected to a tube 30 which is connected to a drip-chamber (not shown) of the piggyback or secondary intravenous set 26. A roller clamp 14a, 30a is typically used on both the primary and piggyback/secondary intravenous sets to control liquid flow rates into the patient.\nThere are several problems associated with the standard techniques employed for intravenous therapy. If the piggyback needle 28 is not securely taped to the Y-injection site 16 and the primary intravenous administration set tubing 12, 18, patients can move their arms, or roll over in bed accidentally pulling the needle 28 completely out of the septum 17a on the Y-injection site 16. If this occurs, the self-sealing latex septum 17a will close off automatically and not allow liquid or contamination to enter the primary intravenous administration set 10. The primary intravenous solution in the bottle or bag 12 will continue to flow into the patient, but, the necessary supplemental pharmaceutical solution from the piggyback or secondary bottle or bag 26 will no longer flow into the patient via the Y-injection site 16. The consequences to the patient for this situation can lead to serious complications and even death if the condition is not noticed by a clinician immediately. Even if the clinician notices the detachment of the needle 28 from the Y-injection site septum 17a immediately, the needle 28 is now contaminated with pathogens and should never be introduced back into the septum 17a. A new sterile, piggyback/secondary intravenous set must be opened, primed, and a new hollow-bore needle reinserted back into the septum on the Y-injection site. Additionally, whether the needle 28 is accidentally detached, or, the clinician removes the needle 28 from the Y-injection site septum once the supplemental pharmaceutical therapy is completed for the patient, the exposed needle 28 is contaminated with pathogens and must be safely disposed by the clinicians without accidentally sticking themselves.\nTo avoid accidental removal, the needle of the piggyback or secondary intravenous set may be taped to the Y-injection site and extension port. When this occurs, the needle may be secured from detachment, but the needle cannot be easily and safely removed by the clinician when the supplemental pharmaceutical therapy is completed, thereby creating a higher incidence for an accidental needle stick injury. Also, because clinicians hold the Y-injection site with one hand while the other hand is used to insert the needle into the Y-injection site septum, the clinicians may accidentally stick the needle directly into their hands holding the Y-injection site, or stick the needle completely through the Y-injection site wall into their hands.\nThe above description and problems associated with conventional continuous and supplemental intravenous therapy through a Y-injection site is similar to the problems associated with intermittent intravenous therapy using a \"Heparin Lock\" injection port 40 (FIGS. 2 and 4). A heparin lock injection port 40 is either connected directly to the vein-access device 20, or attached to a short catheter extension tubing set 42 typically with microbore tubing which is attached to the vein-access device as shown in FIG. 2. The heparin lock has a self-sealing septum port 44 which is similar to the septum port 17 described above. A conventional intermittent intravenous therapy could utilize a short-term primary intravenous administration set 26 with a hollow-bore needle 28 attached to the distal end of a tube 30. The needle would be inserted to the self-sealing septum found on standard heparin lock injection port 40. Another means of introducing supplemental intermittent pharmaceuticals to a patient is to perform an intravenous push utilizing a syringe with a hollow-bore needle attached. The drug is pushed into the patient through the heparin lock injection port 40. Once dispensed, the syringe/contaminated needle is removed from the self-sealing septum 44 on the heparin lock injection port 40.\nAs set out above, the common medical techniques for delivering supplemental liquid fluids to the patient necessitates the use of a hollow-bore needle. The needle is either attached to a secondary intravenous set or a syringe, and is inserted through the self-sealing rubber stopper on the heparin lock injection port or the Y-injection port that is integrated into the primary intravenous administration set. Typically, the needle is secured to the injection port only with tape. The needle can detach from the injection port resulting in a serious or fatal interruption of the flow of the intravenous solutions to the patient. Moreover, the exposed needle can easily be contaminated by contact with non-sterile objects. Sound aseptic techniques must be practiced by the healthcare professional in order to ensure that the sterile needle does not become contaminated and cause a nosocomial infection to the patient.\nSince the discovery of the HIV virus that causes AIDS in the mid-1980's, a major concern among healthcare workers practicing the standard methods of delivering intravenous therapy is accidental needle sticks with a contaminated needle. When a needle is removed from an injection port, it may be contaminated with the patient's blood. The contaminated needle must be carefully disposed in a sharps container. While handling the needle during removal and disposal, clinicians may, and often do, inadvertently stick themselves. Among all of the needled medical devices used in healthcare facilities, contaminated intravenous needles are responsible for the most accidental needlestick injuries. When a needlestick injury occurs, the clinician must stop work and have a blood test performed. Since a needlestick injury can result in fatal disease, the injured clinician will also experience tremendous emotional trauma.\nThere is a wealth of prior art concerned with the problem of accidental needlesticks. Needleless valves are known in the art of intravenous administration. My prior U.S. Pat. No. 5,395,348 discloses a needleless intravenous quick connect/disconnect assembly which includes a first cylindrical member having a female luer connector extending from one end and providing a fluid path to the interior of the first cylindrical member, wherein a first seat is formed in the first cylindrical member at the junction of the female luer and the cylindrical opening; a second cylindrical member having a male-luer connector extending from one end and providing a fluid path to the interior of the second cylindrical member, wherein a second seat is formed in the second cylindrical member at the junction of the male-luer and the cylindrical opening; and a valve member having an integral stem, sealing surface, and resilient one-piece body. In the assembled state, the stem of the resilient valve member extends into the female luer from the inside of the first cylindrical member, and the resilient body of the valve member engages the second seat in the second cylindrical member and biases the sealing surface against the first seat in the first cylindrical member, thereby blocking fluid communication between the female-luer and the interior of the first cylindrical member. The integral valve member is a stepped diameter cylinder having a hollow body and solid stem. The body has a larger diameter than the stem and the step between the body and the stem is frustroconical wherein the sealing surface is formed. The end of the stem is provided with a diametrical slot. An annular fluid passage exists between the valve member body and the cylindrical members. The second seat is preferably formed from a plurality of radially arranged vanes which enter the hollow body of the valve member and flare outward from it. Spaces between the vanes provide a fluid passage from the male-luer to the annular fluid passage surrounding the valve member body. When a male-luer is inserted into the female luer connector, the valve member stem is pressed partially into the first cylindrical member thereby moving the sealing surface away from the first seat and also moving either the diametrical slot in the valve member stem at least partially inside the first cylindrical member. Fluid is free to pass from the male luer through the diametrical slot in the valve member stem, out into the first cylindrical member, around the sealing ring, into the annular fluid passage surrounding the valve member body, through the spaces between the radial vanes, and through the male luer connector in the second cylindrical member. When a male luer is withdrawn from the female luer connector, the resilient body of the valve member biases the valve stem back into the female-luer connector and thereby moves the sealing surface against the first seat, closing the valve.\nDespite the many improvements in the connector disclosed in my U.S. Pat. No. 5,395,348, the valve member inhibits optimal performance of a needleless valve during intravenous and blood collection applications in several ways. Due to the variations of male luer connectors around the world, the one-piece resilient valve stem member can restrict fluid-flow, and the variation of fluid flow from one connector to another can resultingly be unacceptable. Further, the hollow body of the one-piece valve member does not permit optimal blood flushing of the valve, as the hollow body has an open end or concavity in which fluid may accumulate without passing along the flow-stream through the connector.\nThe device disclosed in the above-referenced parent application overcomes many of the shortcomings of the devices of my earlier patent. In the parent application, medical intravenous administration connectors are disclosed which include a first coupling member having a female luer, a valve member having a substantially rigid stem and a substantially resilient body with a sealing surface, and a second coupling member having a fluid coupling extending from one end and an internal valve member support. The coupling members are structured to couple to each other with the valve member being biased to a closed position. When assembled, the valve stem extends into the female luer, and the valve body biases the sealing surface against an annular ring in the first coupling member thereby blocking fluid communication. Preferably, vanes are provided in the second coupling member on which the resilient body of the valve sits, with the vanes acting as a centering mechanism for the valve. The valve may be opened for fluid flow through the assembly by coupling a male luer to the female luer of the assembly, or by pressure actuation. The connectors are easy to prime, limit priming volume, and are arranged without dead-spaces in which blood can be trapped. However, even with these advantages, it is still beneficial to improve the connectors by simplifying the manufacture of the connectors and enhancing their fluid flow performance."}
{"text": "The identification and selection of potentially active drugs for use in treatment of diseases, (for example, cancer) has largely been based on the screening of large numbers of compounds and rational drug design (when the molecular structure of the target such as tumor is known). Once compounds have been chosen based on selectivity to the target and desired functional effect, the ultimate test for advancement of a compound to clinical testing is to show its safety and its efficacy in multiple animal models.\nThere are several current methods and devices that are used to assess the growth of abnormal external features such as tumors, lesions, sores, port stains and wounds in such animal models. The simplest is to measure a feature in two dimensions, length and width, with the aid of a mechanical or electronic caliper that is manually operated. An approximate volume formula is then used to calculate the area or the volume of the object by assuming an ellipsoid shape of the mass. This is the most commonly used method, since it is simple and the calipers cost very little.\nUnfortunately, this approach has obvious pitfalls related to the instrument and methodology used to calculate volume as well as operator-related errors. First, a major source of inaccuracy or variability is caused by the operator's judgment on the choice of the ellipsoid axes. For example, in making a measurement, a different operator may select a different choice of ellipsoid axes. Second, this approach is fairly time consuming. In the case of cancer studies, for example, a single study may involve hundreds of animals, and because the tumor volume measurements have to be done at least twice a week and take a given amount of investigator time per measurement, there is a limit in the number of studies that a single investigator can handle at once. Therefore, a different approach is needed to eliminate the inaccuracies of the methodology and instrumentation used and the subjectivity introduced by the operator.\nOther mechanical volume measuring devices come closer to being an automated volume measurement system. These mechanical devices include a cylindrical chamber filled with elongated pins or filaments which, when pressed over a tumor or any other mass, convert the shape of the mass into volume information. Such devices require direct contact to features such as tumors, however.\nThe most sophisticated methods for measuring a mass include Computer Aided Tomography (CAT) using X-rays, Magnetic Resonance Imaging (MRI) using magnetic scans, Positron Emission Tomography (PET) using electrons and Single Photon Emission Computer Tomography (SPECT) using gamma rays. Such methods are capable of obtaining more accurate size and volume information about both external and internal features such as tumors and organs. The use of such methods is not feasible in many cases, however, due to long and costly preparation times, the unwanted effects of the contrast agents and radioactive compounds and the necessity to anesthetize animals during measurement.\nNotwithstanding the availability of the foregoing techniques, there is no method or device for non-contact and fast measurement of the volume of surface features of small animals that can be used in, for example, pre-clinical cancer and inflammation studies. All the methods discussed above suffer from a number of disadvantages. For example, manual caliper measurements are prone to significant subjective operator errors; CAT, PET, MRI and SPECT methods are unacceptably time consuming and expensive; and caliper and cylindrical chamber methods require direct contact with the object under measurement, which is a potential source of feature deformation and contamination that may negatively effect the measurements being made.\nAs a result there is a need for methods and systems that accurately and quickly measure the volume/size of surface features such as wounds on the surface of animals, plants and humans."}
{"text": "In 1755 LeRoy passed the discharge of a Leyden jar through the orbit of a man who was blind from cataract and the patient saw “flames passing rapidly downwards.” Ever since, there has been a fascination with electrically elicited visual perception. The general concept of electrical stimulation of retinal cells to produce these flashes of light or phosphenes has been known for quite some time. Based on these general principles, some early attempts at devising prostheses for aiding the visually impaired have included attaching electrodes to the head or eyelids of patients. While some of these early attempts met with some limited success, these early prosthetic devices were large, bulky and could not produce adequate simulated vision to truly aid the visually impaired.\nIn the early 1930's, Foerster investigated the effect of electrically stimulating the exposed occipital pole of one cerebral hemisphere. He found that, when a point at the extreme occipital pole was stimulated, the patient perceived a small spot of light directly in front and motionless (a phosphene). Subsequently, Brindley and Lewin (1968) thoroughly studied electrical stimulation of the human occipital (visual) cortex. By varying the stimulation parameters, these investigators described in detail the location of the phosphenes produced relative to the specific region of the occipital cortex stimulated. These experiments demonstrated: (1) the consistent shape and position of phosphenes; (2) that increased stimulation pulse duration made phosphenes brighter; and (3) that there was no detectable interaction between neighboring electrodes which were as close as 2.4 mm apart.\nAs intraocular surgical techniques have advanced, it has become possible to apply stimulation on small groups and even on individual retinal cells to generate focused phosphenes through devices implanted within the eye itself. This has sparked renewed interest in developing methods and apparatus to aid the visually impaired. Specifically, great effort has been expended in the area of intraocular retinal prosthetic devices in an effort to restore vision in cases where blindness is caused by photoreceptor degenerative retinal diseases; such as retinitis pigmentosa and age related macular degeneration which affect millions of people worldwide.\nNeural tissue can be artificially stimulated and activated by prosthetic devices that pass pulses of electrical current through electrodes on such a device. The passage of current causes changes in electrical potentials across visual neuronal membranes, which can initiate visual neuron action potentials, which are the means of information transfer in the nervous system.\nBased on this mechanism, it is possible to input information into the nervous system by coding the sensory information as a sequence of electrical pulses which are relayed to the nervous system via the prosthetic device. In this way, it is possible to provide artificial sensations including vision.\nOne typical application of neural tissue stimulation is in the rehabilitation of the blind. Some forms of blindness involve selective loss of the light sensitive transducers of the retina. Other retinal neurons remain viable, however, and may be activated in the manner described above by placement of a prosthetic electrode device on the inner (toward the vitreous) retinal surface (epiretinal). This placement must be mechanically stable, minimize the distance between the device electrodes and the visual neurons, control the electronic field distribution and avoid undue compression of the visual neurons.\nIn 1986, Bullara (U.S. Pat. No. 4,573,481) patented an electrode assembly for surgical implantation on a nerve. The matrix was silicone with embedded iridium electrodes. The assembly fit around a nerve to stimulate it.\nDawson and Radtke stimulated cat's retina by direct electrical stimulation of the retinal ganglion cell layer. These experimenters placed nine and then fourteen electrodes upon the inner retinal layer (i.e., primarily the ganglion cell layer) of two cats. Their experiments suggested that electrical stimulation of the retina with 30 to 100 μA current resulted in visual cortical responses. These experiments were carried out with needle-shaped electrodes that penetrated the surface of the retina (see also U.S. Pat. No. 4,628,933 to Michelson).\nThe Michelson '933 apparatus includes an array of photosensitive devices on its surface that are connected to a plurality of electrodes positioned on the opposite surface of the device to stimulate the retina. These electrodes are disposed to form an array similar to a “bed of nails” having conductors which impinge directly on the retina to stimulate the retinal cells. U.S. Pat. No. 4,837,049 to Byers describes spike electrodes for neural stimulation. Each spike electrode pierces neural tissue for better electrical contact. U.S. Pat. No. 5,215,088 to Norman describes an array of spike electrodes for cortical stimulation. Each spike pierces cortical tissue for better electrical contact.\nThe art of implanting an intraocular prosthetic device to electrically stimulate the retina was advanced with the introduction of retinal tacks in retinal surgery. De Juan, et al. at Duke University Eye Center inserted retinal tacks into retinas in an effort to reattach retinas that had detached from the underlying choroid, which is the source of blood supply for the outer retina and thus the photoreceptors. See, e.g., E. de Juan, et al., 99 Am. J. Ophthalmol. 272 (1985). These retinal tacks have proved to be biocompatible and remain embedded in the retina, and choroid/sclera, effectively pinning the retina against the choroid and the posterior aspects of the globe. Retinal tacks are one way to attach a retinal electrode array to the retina. U.S. Pat. No. 5,109,844 to de Juan describes a flat electrode array placed against the retina for visual stimulation. U.S. Pat. No. 5,935,155 to Humayun describes a retinal prosthesis for use with the flat retinal array described in de Juan."}
{"text": "Cranes and excavators, as well as various other types of agricultural, industrial and construction machines, have heretofore been provided with wheels which are connected at fixed locations to the chassis and with outriggers which may be provided to support the chassis or vehicle body in various attitudes and to various degrees upon the ground. The outriggers may be formed with legs which are telescopingly extensible so that a pad or the like can rest upon the ground and can support the vehicle body which may be lifted to elevate the wheels above the ground.\nIt is also known to provide wheel-carrying assemblies which are provided on rigid traveling legs or legs which are hinged at fixed locations to the chassis to cooperate with a vehicle body which is capable of rolling and/or stepping movement along the ground.\nFor example, the vehicle body may have arms which are themselves hinged to a central support member so that the vehicle body may be raised into a portal or gantry configuration or lowered into a compact configuration with the center of the body close to the ground.\nSuch traveling outriggers enable the desired attitude of the vehicle to be maintained when the vehicle is used upon irregular terrain.\nOne of the problems encountered with earlier systems in which rolling and/or stepping travel is required, is that the usual wheel suspension or support does not permit rolling travel or travel on steep, banked or stepped terrain to a satisfactory degree."}
{"text": "The present invention relates to a differential apparatus to which driving force is to be cut off in a part-time four(4)-wheel drive car when in two(2)-wheel drive. More particularly, the invention relates to a differential apparatus having a mechanism for reducing frictional loss of the apparatus when not driven.\nIn prior art part-time four-wheel drive cars of this kind, a differential apparatus such as shown in FIG. 1 is used. Such an apparatus is disclosed in Japanese Utility Model Publication (kokai) No. 63-193929.\nIn FIG. 1, assembled so as to rotate in a first casing 401 relatively is a second casing 403 in which a pinion shaft 405, a pinion gear 407 and left and right side gears 409, 411 etc. are arranged. A sleeve 413 is arranged so as to slip on the outer periphery of the first casing 401. With the slipping movement of the sleeve 413, both casings 401, 403 can be either operatively associated with, or separated from, each other.\nWhen it is desired to switch from four-wheel drive into two-wheel drive, due to a driver's manipulation of a controller in a vehicle cabin, the sleeve 413 is shifted to the left side of FIG. 1 by a switching unit 417 provided on a differential carrier 415. Consequently, the association between the first casing 401 and the second casing 403 is released, so that the casings can rotate relatively to each other.\nIn this way, while the second casing 403, the pinion shaft 405, the pinion gear 407 and the left and right side gears 409, 411 etc. race independent of tire rotation due to the vehicle's traveling, the surrounding first casing 401 and a ring gear 419 etc. do not race since the power transmission is cut off by switching of a transfer unit. Thus, racing loss and occurrence of noise can be reduced.\nIn the above-mentioned differential apparatus, however, since the switching unit 417 for switching between 4-wheel/2-wheel driving operation is disposed outside the differential carrier 415, it has been unavoidable that the differential carrier 415 is subjected to detailed remodeling to install a mechanism for sliding the sleeve 413 in the carrier 415.\nFurthermore, the known art has to ensure a space allowing the axial sliding of the sleeve 413 in the first casing 401. For this reason, the first casing 401 is apt to have a one-sided weight imbalance and deviated configuration. Thus, an ill balanced configuration is caused in the differential apparatus."}
{"text": "In medical practice, syringe pumps are used for drugs which need high accuracy and have a short half-life in the body. In operating theatres and intensive care units, syringe pumps are mounted in stacks of in particular six to eight pumps that, however, require a lot of space and allow only low visibility and recognition of the indication of a drug provided on the syringe. Apart from occupation of large space syringe pumps have more problems, since a nominal trumpet curve and, thus, a constancy index cannot be achieved in practice due to the fact that the plunger usually made of rubber sticks to the walls of the syringe and therefore advances in pulses rather than continuously. Further, syringe pumps also have low sensitivity in occlusion pressure reading, that becomes a problem in neonatal infusions and recently with the use of wearable bolus large volume injectors. Syringe pumps are extensively used mainly in Europe, where about 40% of the beds in each hospital are equipped with a syringe pump, and for insulin infusions and most of immunoglobulin and Parkinson's disease infusions. It has been proposed to replace insulin syringe pumps by diaphragm pumps that, however, cannot be realized in practice since insulin crystallizes and renders active and passive valves of the mechanism leak.\nPre-filled syringes are part of a growing pharmaceutical delivery sector and work well for injections, but are problematic for longer term infusions, since they become bulky especially for newer biological drugs and have a volume limit of about 60 to 100 ml. So, the pumps become bulky as the needed infusion volume increases.\nSyringe pumps are used because prior art peristaltic pumps had a low accuracy and causes a pulsatile flow and, hence, not a linear flow per infusion cycle, wherein during a part of the cycle there is no infusion but sometimes even a backflow, so that their constancy index is high for short half-life drugs. Short-term accuracy can be expressed by the concept of constancy index. This is the shortest period during a steady-state operation of a pump over which a measurement of output consistently falls within +/−10% of the mean rate. These data are derived from flow tests performed over 24 hours at 1 ml/h, wherein the flow is recorded at 30 seconds intervals during the final 18 hours period and the average rate is compared with the flow over each short period.\nPeristaltic pumps comprise a housing and a compressible tubing arranged within the housing. Basically, there are two different embodiments of the peristaltic pump, wherein in the one embodiment the tubing is arranged along a straight track, whereas in the other embodiment the tubing is formed as a loop resulting in a more economical design with a smaller physical size and less producing costs. The former embodiment which is called a linear peristaltic pump is mainly used nowadays, while the present invention deals with the latter embodiment. The tubing is to be filled with a fluid to be delivered from its inlet to its outlet. The fluid is caused to move through the tubing by engagement elements, typically in the form of rollers driven by rotary means such an electric motor or a mechanically driven shaft. The engagement elements cause an occlusion of the tubing by squeezing it against a wall or track within the housing so that the fluid is forced through the tubing due to the movement of the engagement elements. The use of a peristaltic pump is advantageous in so far as the fluid does not come into contact with the operating environment, which renders the peristaltic pump suitable for medical applications like infusion of drugs where it is important to avoid contact of the fluid with the environment. Further, the mechanical components of a peristaltic pump do not come into contact with the fluid. So, the components of a peristaltic pump remain free from contamination by the fluid. As a result, a peristaltic pump is easy to clean and to sterilize because the tubing can be simply discarded after use, and a new tubing can be provided for the next use.\nHowever, a disadvantage is that it is difficult with a peristaltic pump to achieve a constant or pulseless flow of the fluid through the tubing. Pulses are created when the engagement elements disengage from the tubing and, therefore, the occlusion is removed with the result of that a void is generated in the disengagement region of the tubing. Namely, in this region the tubing returns to its normal round shape resulting in an increase of the volume which are filled by fluid from the outlet of the tubing. This leads to a reduction of the flow rate of the fluid at the outlet of the tubing for the duration of the pulse.\nIn other words, the pulsatile behavior of a rotary peristaltic pump results from the alternation of a forerunner or leading engagement element by the next follower or trailing engagement element. V is the volume which is encapsulated between two neighboring engagement elements. In case of a rotary peristaltic pump, with φ defining an angular position so that the angular position of the trailing engagement element is φ1 and the angular position of the leading engagement element is φ2, the volume V encapsulated between both these engagement elements extends along an angular distance which is defined by the difference between both the aforementioned angular positions φ2 and φ1, i.e. Δφ=φ2−φ1. +ΔV/Δφ represents an increase of the volume defining a so-called frontwave which is displaced in front of each engagement element and is advanced by it. −ΔV/Δφ represents a decrease of the volume defining a so-called depression which arises behind each engagement element. The enclosed volume V between two neighboring engagement elements squeezing the tubing with unchanged distance between them is constant so that the fluid is just transported, if +ΔV/Δφ and −ΔV/Δφ are the same so thatV+ΔV/Δφ−ΔV/Δφ=V, resulting in that the pressure remains constant as well.\nIf otherwise an upstream portion of the tubing is larger than a downstream portion so that it isΔV2/Δφ>ΔV1/Δφ andV+ΔV2/Δφ−ΔV1/Δφ=V+Vdifference,wherein ΔV1/Δφ represents an increase of volume defining a front wave in front of the forerunner or leading engagement element and ΔV2/Δφ represents an increase of volume defining a front wave in front of the next follower or trailing engagement element, the pressure is increased by elastic tubing forces of the inflated portion due to increase of volume.\nAt the moment the leading engagement element stops squeezing the tubing and a front/back communication through a thin film of fluid is established under it, the frontwave +ΔV/Δφ suddenly disappears (so that it cannot push fluid anymore) and is replaced by the frontwave in front of the trailing engagement element which now takes over relay of infusion. Also −ΔV/Δφ disappears behind it. But due to the disengagement of the leading engagement element a new additional difference volume ΔVd/Δφ is built up and continues to be present until the disengagement of the leading engagement element is fully completed. A void creating the aforementioned new additional volume difference ΔVd arises, as the resilient tubing becomes asymptotic or starts with a larger diameter disengagement curvature at this point, resulting in a creation of a negative pulse ΔV/Δφ−ΔVd/Δφ in the flow graph (V,φ) (where ΔVd depends on the geometry of the engagement element and the disengagement curvature Δr/Δφ of the housing and r is the radius of the disengagement curvature increasing so as to let the engagement element disengage).\nThe flow rate is defined as ΔV/Δt. By multiplying the nominator and the denominator each with Δφ, the above equation can be written asflow rate=(ΔV/Δφ)·(Δφ/Δt).\nThis makes evident that for a constant rotational movement Δφ/Δt, if ΔV is constant at the vicinity of an engaging element at any location along its movement path, the flow rate will be constant without any variation or pulse.\nThe lack of a constant fluid caused by pulses in the tubing renders a peristaltic pump unsuitable for certain precision applications. E.g., in applications where a small volume of fluid is required, such as where less than a complete revolution of the rotor is used, the effect of the pulses are particularly disadvantageous.\nU.S. Pat. No. 5,533,878 A discloses a squeeze type pump wherein the resilient tubing has a larger diameter at the start of an infusion cycle than at its end.\nU.S. Pat. No. 7,654,127 B2 discloses a peristaltic pump with increased dimensions in an upstream portion of the tubing so as to increase pressure inside the tubing before disengagement of a roller squeezing the tubing. However, the pressure is suddenly released when the leading roller is going to be disengaged from the tubing. Therefore, a perturbation of flow happens first with a sudden increase of flow and second with a decrease of it at same quantity due to disengagement of the roller from the tubing after some rotational degrees.\nU.S. Pat. No. 3,826,593 A proposes the provision of a cam compressing the tubing at another point at the same time the engagement roller is disengaged from the tubing in a peristaltic pump. However, this solution requires higher costs and a larger number of assembly parts and is therefore not suitable for a disposable pumping mechanism.\nU.S. Pat. No. 5,470,211 A teaches a roller pump with the provision of a controlled curvature at the inlet and the outlet of the tubing.\nIt is an object of the present invention to provide a continuously operating pulseless and accurate peristaltic pump having a small size so that it can replace many bulky syringe pumps in limited space."}
{"text": "The invention relates to a turbomachine, as described and shown below.\nIn order to diminish structural component temperatures in the region of the final compressor stage and also a rotor shaft segment adjoining the compressor stage, cooling air is preferably blown in from an air distribution chamber between an outlet diffuser of the compressor and an associated combustion chamber, instead of using cooling air from the boundary layers in the flow conduit of the compressor in the region of the inner shrouds. The supply of cooling air into a cavity from the region in front of the combustion chamber occurs conventionally through a plurality of successively arranged intervening chambers of the outlet diffuser, which are in fluid connection with one another. The drawback of this cooling air supply, however, is that, owing to the plurality of intervening chambers, the cooling air undergoes a significant loss of pressure. If, in addition, radial supply lines, such as air bleed or oil lines traverse an intervening chamber, the pressure loss to which the cooling air is subject is augmented. In addition, leakage can occur in the region of inlet and outlet points of the supply lines due to the intervening chambers."}
{"text": "The present invention relates to a closeable opening device for a sealed packaging, which contains a flowable medium, wherein the device comprises a spout fitting having a pipe section and flange, said pipe section being furnished with an internal and an external thread, as well as to a cylindrical cutting element, the lower edge of which is equipped with one or a plurality of teeth and the internal jacket surface of which is furnished with driving dogs protruding in the radial direction that interact with driving dogs in a screw cap comprising a jacket wall and cover surface, wherein the spout fitting and the cylindrical cutting element are produced as a one-piece finished good connected to one another via predetermined breaking points.\nA closeable opening device of the aforementioned type is taught by the European Patent Reference EP-A 1088764. In this case, a principle is used, which is common in closure technology, wherein multi-part closures are produced so that two parts are disposed on top of one another and are injection molded in one piece to be connected to one another via predetermined breaking points. In this connection, the expenditure for tools is on the one hand reduced and on the other hand the assembly is simplified because the parts connected to one another via predetermined breaking points are already disposed in their accurately aligned relative position to one another and have thus to be merely pushed together. This technology, which is sufficiently known per se, has also been used previously for the closures of interest here and is, for example, also taught by the European Patent Reference EP-A 1084060.\nBoth documents disclose a semi-finished good, which displays a one-piece production of a spout fitting including a flange together with a cylindrical cutting element. The cutting element according to the European Patent Reference EP-A 1084060 does not have a planar surface at the upper or lower edge, which is suitable for being connected in a revolving manner to an upper or lower edge of the spout fitting. A solution is thus necessarily disclosed in which the cutting element is produced in an intermediate position within the spout fitting as a semi-finished good connected to the lower edge of said spout fitting via predetermined breaking points.\nIn both solutions, the preassembly, during which the cutting element has to be inserted into the spout fitting, is not as trivial as it may appear. Sealed packagings made of laminates or plastic films, onto which the opening devices of interest here are mounted, are or contain a plastic film which is difficult to sever because the material is very tough and requires that the teeth correspondingly comprise a sufficient solidity. Nevertheless, said teeth cannot simply be designed arbitrarily robust because they also must be capable of piercing and cutting the material. Correspondingly said teeth must not be negatively impacted during assembly. This problem is solved by WIPO patent application WO 2007/051328. The solution consists of holding the cutting element and the spout fitting together such that a simple axial pressure can be applied to the upper edge, which runs flatly in a perpendicular plane with respect to the rotational axis. If this assembly is carried out, the screw cap must now be fitted so as to be correctly aligned in the axial direction. This is carried out to this point as a result of the screw cap being fitted practically in an arbitrary angular position with respect to the perforator and the screw cap being subsequently pressed with the internal thread thereof in a ratchet-like manner over the external thread on the spout fitting. This is not a problem per se with regard to the assembly. A problem has however become apparent in that the driving dogs on the screw cap and the corresponding driving dogs on the cutting element can now practically assume an arbitrary angular position with respect to each other. Because the tamper-evident band has to first be detached during the initial opening and in so doing a relatively large torque is required, the user tends to apply said torque in a sudden and jerky manner. This can now result in the driving dogs of the screw cap being able to jump over the driving dogs of the spout fitting.\nIn addition, the greatest torque is required when the teeth have to open the film of the sealed packaging. At this time, it is again possible that the driving dogs jump over one another. If this happens, a complete opening of the film is often no longer assured when the packaging is being initially opened.\nIn order to solve the first problem, assembly lines have partially been built, which have turned the screw cap into an exact angular position with respect to the cutting element and only then push the two parts together so that in the assembled state the driving dogs are already brought correctly to the position where said dogs approximately abut against each other. This has required corresponding positioning means on the screw cap as well as on the semi-finished product, which can be mechanically or optically scanned. This process is, however, associated with substantially higher assembly costs.\nWith regard to the second problem, a practical solution thereto has so far not been found."}
{"text": "One embodiment of a prior power converter 100 that may operate in transition mode is schematically depicted in FIG. 1. In this regard, transition mode may be defined as a boundary mode between a discontinuous operating mode and continuous operating mode. The power converter 100 of FIG. 1 includes a transformer (T1) 110, two MOSFETS (M1, M2) 120 and 122, four diodes (D2, D4, D5, D6) 130-136, three capacitors (C3, C4, C5) 140-144, and a resistor (R5) 152. The power converter 100 is shown connected to a direct current voltage source 160. The transformer 110 includes primary and secondary windings (T1-1, T1-2) 110A and 110B. The MOSFETS (M1, M2) 120 and 122 may be referred to herein as the primary MOSFET (M1) 120 and the auxiliary MOSFET (M2) 122. In FIG. 1, parasitic capacitances and inherent body diode characteristics of the primary and auxiliary MOSFETS (M1, M2) 120 and 122 are not illustrated. Diodes (D4, D5, D6) 132-136 and capacitor (C5) 144 together comprise driving circuitry for auxiliary MOSFET (M2) 122 to control switching of auxiliary MOSFET (M2) 122 between an on and an off state, preferably under substantially Zero Voltage Switching (ZVS) conditions. Rectifier diode (D2) 130 and smoothing capacitor (C3) 140 together comprise secondary side rectification means.\nAs depicted in FIG. 1, the prior transition mode power converter 100 technology has external capacitor (C5) 144 and switch diode (D4) 132 arranged in series with the gate input capacitance of auxiliary MOSFET (M2) 122. This serial combination (C5, D4, M2 gate input capacitance) is in parallel with snubber capacitor (C4) 142 to divide the high voltage on snubber capacitor (C4) 142 into a certain level of voltage to safely drive the gate of auxiliary MOSFET (M2) 122 to an on state. MOSFET (M2) 122 can be turned off when the voltage across snubber capacitor (C4) 142 drops to zero using the other two switch diodes (D5 and D6) 134 and 136.\nThe divider circuit used in the prior transition mode power converter of FIG. 1 is subject to the tolerance of auxiliary MOSFET (M2)'s 122 gate input capacitance, which usually is not clearly defined and can vary by drain-source voltage as well as lot by lot and vendor by vendor. Additionally, the voltage on snubber capacitor (C4) 142 can vary significantly with the output load and the input line voltage as well. Thus, the auxiliary MOSFET (M2)'s 122 gate voltage level can have significant variation, especially if the output load range is large. Considerations such as these can make the prior power converter 100 of FIG. 1 less desirable for higher power applications."}
{"text": "In today's society, users are increasingly utilizing network and other service providers to gain access to the Internet, access software services, request and receive various types of content, access software applications, and perform a variety of other tasks and functions. In order to handle large volumes of requests from users requesting various types of content, services, and applications, content providers often employ the use of content delivery networks and systems to serve the content and provide the services to the users in a more efficient and effective manner. Content delivery networks are particularly useful when high performance, reliability, scalability, and cost savings associated with distributing the content and providing the services are important to the content providers that provide the content, services, and software applications. Additionally, content delivery networks are useful in optimizing the delivery of video content to a multitude of users and organizations based on their bandwidth constraints.\nWhile current implementations of content delivery networks provide substantial benefits to users and content providers, there is still significant room to enhance current methodologies and technologies for augmenting the content that is delivered by content delivery networks and improving the content delivery experience for users. For example, currently, on a large display device, various types of small or micro-movements may be easily observed by one or more users. However, on devices with limited display capabilities or on smaller devices, such small or micro-movements may not be readily observable. Additionally, for certain types of video content and environments, there may be excessive movement that may interfere with the viewing experience. For example, when a user is on a train, the movement of the train travelling on a set of train tracks may interfere with the user's viewing of media content. As a result, content delivery networks and the various methods of utilizing content delivery networks to deliver media content may be modified so as to provide enhanced quality-of-service for users. Such enhancements and improvements to content delivery networks may provide for improved customer satisfaction and increased revenue for content providers and service providers."}
{"text": "1. Field of the Invention\nThe present invention relates to a toner for developing an electrostatic image in dry process to form images in electrophotography, electrostatic recording, electrostatic printing, and the like. The present invention also relates to a fine powdery titanium oxide, particularly to hydrophobic fine powdery titanium oxide suitable as an external additive to the toner.\n2. Related Background Art\nMethods for developing an electrostatic image are well known in which an electrostatic image is formed on a photoconductive surface by electrostatic means and subsequently the image is developed with a toner. Various methods have been disclosed, for example, in U.S. Pat. No. 2,297,691, JP-B-42-23910, and JP-B-4324748 (\"JP-B\" herein means \"examined Japanese patent publication). In these methods, generally, an electrostatic latent image is formed on a photosensitive member employing a photoconductive substance and thereon a toner image is formed corresponding to the electrostatic image with the toner adhering to the latent image. Then, the toner image is transferred onto a surface of a desired image supporting medium (recording medium) such as a paper sheet, and the transferred image is fixed by heating, pressing, heat-pressing, or solvent-vapor exposure to obtain a copy. In the method involving a step of transferring the toner image, a cleaning step is usually provided to remove remaining toner as required.\nThe known developing methods to visualize an electrostatic image with a toner include powder cloud development described in U.S. Pat. No. 2,221,776, cascade development described in U.S. Pat. No. 2,618,552, magnetic brush development described in U.S. Pat. No. 2,874,063, and electroconductive magnetic toner development described in U.S. Pat. No. 3,909,258.\nThe toner used in such development methods is generally prepared by mixing a thermoplastic resin with a colorant for dispersion, and pulverizing it. Polystyrene resin is generally used as the thermoplastic resin. Polyester resin, epoxy resin, acrylic resin, urethane resin, or the like is also useful therefor. As the colorant, carbon black is widely used. In a magnetic toner, iron oxide type black magnetic powder is widely used. When a two-component developer is used, the toner is usually mixed with a particulate carrier such as glass beads, iron powder, magnetic ferrite powder, or resin-coated particles thereof.\nThe toner image on the final image bearing member like a paper sheet is fixed thereon by heating and/or pressing. Thermal fixation is frequently employed therefor.\nIn recent years, copying machines and printers have been rapidly developing from mono-color image to full-color image, and two-color copying machines and full-color copying machines are now practically used. For example, studies on color reproduction, or tone reproduction are reported in \"Denshi Shashin Gakkaishi (Journal of Electrophotography Society)\": Vol. 22, No.1 (1983), and Vol. 25, No.1, page 52 (1986).\nFurther improvement of the present full-color electrophotographic image formation method is required, since the people are accustomed to artificial color images like television images, color photographs and color prints more beautiful than actual objects and not compared directly with the actual objects.\nIn full-color electrophotography, all colors of an image are generally reproduced with three kinds of primary color toners, that is, yellow, magenta, and cyan. In this method, an electrostatic latent image is formed on a photoconductive layer using light from the original passing through a color separation filter which transmits only the complementary color of a toner color, and then the image is developed with the toner and transferred onto a recording member. This process is repeated for each color, and the toner images are superposed on one and the same recording member with registration, and the superposed images are fixed in one step to obtain a final full-color image.\nWhen a two-component developer composed of a toner and a carrier is used, the toner particles are electrostatically charged by friction with the carrier particles to a required level in the desired polarity, and the electrostatic image is developed using electrostatic attraction. Therefore, the triboelectrification properties of the toner mainly determined in relation to the carrier are important for obtaining a excellent visible image.\nAccordingly, to obtain satisfying triboelectrification properties, many studies have been made on carrier core materials, carrier core-coating materials, optimum coating amount, toner charge-controlling agents, fluidity-improving agents, improvement of the toner binder, and so forth.\nFor example, an electrification aid like an electrifiable fine particulate material is added to the toner, JP-B-52-32256 and JP-A-56-64352 (\"JP-A\" herein means \"Unexamined Laid-Open Japanese Patent Application\") propose to add to the developer a fine powdery resin having the polarity opposite to the toner; and JP-A-61-160760 discloses addition of a fluorine-containing compound to the developer to obtain stable triboelectrification properties of the toner.\nThe above electrification aid is generally added by adhesion of the electrification aid to the toner particle surface through electrostatic force, Van der Waals force, etc. where a stirring or mixing machine is used for it. In such a method, uniform adhesion of the electrification aid does not readily occur on the toner particle surface, and the additive not adhering to the toner form aggregates. This tendency is more remarkable with the larger electric resistivity and the smaller particle size of the electrification aid, which adversely affects the toner performance. For example, unstable triboelectrification, fluctuation of image density, and fogging tend to occur. There is also a problem that the content of the electrification aid changes so that the initial image quality cannot be maintained during continuous copying.\nIn another method of addition, the electrification aid is added in advance during production of the toner along with a binding resin and a colorant. However, uniform dispersion of the electrification aid is difficult. Further, only the electrification aid on or near the surface actually contributes to the electrification, and the electrification aid or the charge-controlling agent in the interior of the toner particles do not contribute to the electrification properties. Therefore, the amount of the electrification aid to be added or the amount thereof on the toner surface is difficult to control, and the triboelectrification is liable to be unstable even with the toner prepared by such a method.\nIn recent years, for image output of copying machines and printers, finer and higher quality is required. To improve the image quality, it has been attempted to make the toner particle diameter smaller. However, the finer particle diameter of the toner results in increase of the surface area per unit weight, thus increasing the amount of electrification, which tends to cause low image density and deterioration of durability. Additionally, the larger electric charge of the toner increases the adhesion between toner particles, lowering fluidity of the toner and causing problems in stability of toner supply and impartment of triboelectricity to the supplied toner.\nA color toner, which does not contain electroconductive substance like carbon black or magnetic material to leak electric charge, tends to be charged at a higher level. This tendency is remarkable when a polyester binder is employed as the binder for the toner.\nFor color toners, are strongly required properties as set forth in below.\n(1) The fixed toner particles should be almost completely fused so that the individual toner particles are not discriminable to reproduce the color without irregular reflection of light; PA1 (2) The color toner should be transparent not to intercept the underlying color of another color toner layer; and PA1 (3) Each color toner should have well-balanced hue and spectral reflection characteristics, and sufficient color saturation.\nFrom such viewpoints, various binder resins are under investigation to satisfy the above requirements. At the moment, polyester resins are widely used as the binder resin for the color toner. The toner comprising a polyester resin is generally affected by temperature and humidity, electric charge is excess at low humidity, and insufficient at high humidity. Therefore, color toners are desired which have stable charge under various environmental conditions.\nIn JP-A-48-47345, use of powdery metal oxide is disclosed as an abrasive. Metal oxides like titanium oxide are disclosed as the fluidizing agent in JP-A-52-19535 and JP-A-56-128956. Further, a surface treated amorphous titanium oxide powder is disclosed in JP-A-4-337739, JP-A-4-348354, JP-A-4-40467, and JP-A-5-72797 for the purpose of imparting fluidity to toner particles, stabilizing electric charge and preventing the filming of the toner.\nAddition of hydrophobic titanium oxide treated with alkyltrialkoxysilane to the toner is proposed in JP-A-59-52255. The added hydrophobic titanium oxide improves electrophotographic properties of the toner. However, the hydrophobicity may be still insufficient for that purpose since the surface activity of titanium oxide is inherently low in comparison with silica. Further increasing the amount of the treating agent or the treatment period to raise the hydrophobicity is liable to cause agglomeration of titanium oxide particles during the treatment or to cause non-uniformity of the hydrophobicity.\nFine powdery materials are used in various fields and their surfaces are usually treated to meet the requirements in the particular field.\nAs stated above, in the electrophotographic application field, fine powdery silica, fine powdery titanium oxide, or the like is used as the external additive to the toner for electrophotographic development, and the powdery material treated for hydrophobicity is also used to stabilize the electrophotographic properties of the toner in various environmental conditions.\nKnown titanium oxide includes those formed by neutralizing aqueous titanium sulfate solution and firing the formed precipitate; those formed by decomposing and oxidizing titanium chloride at a high temperature; and those formed by hydrolyzing or pyrolyzing titanium alkoxide. These titanium oxide are anatase, rutile, or amorphous in crystal structure. The titanium oxide is treated with a hydrophobicity-imparting agent such as the silane-coupling agent to form hydrophobic titanium oxide. Known titanium oxides, however, are not sufficiently reactive to the hydrophobicity-imparting agent, and therefore development of more reactive fine powdery titanium has been desired."}
{"text": "The present invention relates to a modular clutch construction, particularly to a modular clutch construction wherein the peripheral wall of a clutch cover is marginally fitted at one end onto the cylindrical surface of a flywheel.\nA conventional clutch comprises mainly a clutch cover assembly retaining a pressure plate fixed to a flywheel, and a clutch disc assembly located between the flywheel and the pressure plate of the clutch cover assembly.\nThe clutch cover end extends radially outward in a flange, and the clutch cover assembly thus is mounted to the flywheel by fixing the flange to an end-face thereof. Therein, the outside diameter of the flywheel must be larger than that of the clutch disc assembly by the radial width of the flange.\nJapanese Utility Model Provisional Number No. 1762 of 1975 discloses a clutch cover attachment structure for minimization of flywheel dimensions, wherein the peripheral wall of the clutch cover is circumferentially assembled onto the cylindrical surface of the flywheel. The outer diameter of the flywheel is on the order of that of the clutch cover, since the flange of a conventional clutch cover is eliminated.\nIn the foregoing, bolts fix the clutch cover peripheral wall to the flywheel cylindrical surface. Due to the curvature of the flywheel cylindrical surface, the area of contact in seating the fixing bolts is limited. Consequently, retaining friction between the surfaces is insufficient, such that the bolts tend to loosen and portions in seating contact tend to wear. Moreover, shearing forces on the clutch cover act directly upon the bolts, occasionally breaking them off. A solution to this problem is to increase the number of bolts; however, a larger number of bolts adds complexity to the assembly stage.\nFurthermore, the clutch cover tends to expand radially outward under centrifugal force generated by rotation of the flywheel. As a result, there is a load on the bolts fastening the clutch cover to the flywheel, which tends to loosen them."}
{"text": "Tertiary amides are used as solvents in a variety of industrial chemical processes, and their recycling for reuse in these processes is of significant commercial value. Depending on the type of reaction performed in the tertiary amide, such recycling may require the removal of reaction byproducts not easily separated from it in a cost-effective manner. One example of such byproducts includes carboxylic acids. In some cases, these must be removed from used solvent as part of recycling, and in other cases they must be removed from the solvent while it is still in use to dissolve process intermediates. The latter situation frequently compounds the difficulty of removing the acids, since potentially sensitive intermediates may be present. The commercial preparation of sucralose is an example of a process where carboxylic acids need to be removed from a tertiary amide solvent, both alone and in the presence of sensitive process intermediates.\nSucralose (4,1′,6′-trichloro-4,1′,6′-trideoxygalactosucrose), a high-intensity sweetener made from sucrose, can be used in many food and beverage applications.\nA number of different synthesis routes for the preparation of sucralose have been developed in which the reactive hydroxyl in the 6 position is first blocked with an acyl group to form a sucrose-6-ester. The acyl group is typically acetyl or benzoyl, although others may be used. The sucrose-6-ester is then chlorinated to replace the hydroxyls at the 4, 1′ and 6′ positions to produce 4,1′,6′-trichloro-4,1′,6′-trideoxygalactosucrose 6-ester (referred to herein as sucralose-6-ester), followed by hydrolysis to remove the acyl substituent and thereby produce sucralose. Several synthesis routes for formation of the sucrose-6-ester involve tin-mediated acylation reactions, with illustrative examples being disclosed in U.S. Pat. Nos. 4,950,746; 5,023,329; 5,089,608; 5,034,551; and 5,470,969, all of which are incorporated herein by reference.\nVarious chlorinating agents may be used to chlorinate the sucrose-6-ester, and most commonly a Vilsmeier-type salt such as Arnold's Reagent will be used. One suitable chlorination process is disclosed by Walkup et al. (U.S. Pat. No. 4,980,463), incorporated herein by reference. This process uses a tertiary amide, typically N,N-dimethyl formamide (“DMF”), as the chlorination reaction solvent. After the chlorination is complete, adducts of Arnold's reagent on the base sucrose moiety and excess chlorinating reagent are neutralized (“quenched”) with aqueous base to provide the sucralose-6-ester in an aqueous solution, accompanied by the tertiary amide solvent and salts resulting from reactions of the chlorination reagent. The sucralose-6-ester is then deacylated to produce sucralose. One suitable deacylation process is taught by Navia et al, U.S. Pat. No. 5,498,709, the entire disclosure of which is incorporated herein by reference.\nIn such processes, carboxylic acids need to be removed at various points in the process. Accordingly, facile means of removing these acids are of commercial value."}
{"text": "1. Field of the Invention\nThis invention is in the field of seals for fluids; more particularly, the invention relates to seals for use in rotary pumps, engines, and the like.\n2. Description of the Prior Art\nDevices for handling fluids in various ways are known in the art, and are useful in many applications, such as in compressors, pumps, and external and internal combustion engines.\nIn recent years, rotary engines have been the subject of a great deal of engineering development. The most prominent among these is an internal-combustion engine commonly called the Wankel. This engine comprises a rotary member mounted on eccentric bearings, and operating in a trochoidal chamber. Rotation of the rotary member is kept in phase by the use of phasing gears, one of which is a stationary reaction gear fixed to the end cover plates mounted concentrically with the mainshaft.\nBecause of the simplicity of construction, the Wankel-type engine offers many advantages such as long life, low vibration, and high power-to-weight ratio. However, there are some problems with this engine which have persisted despite considerable effort spent on their solution. The principal one among these is that of seals in general, and apex and corner seals in particular.\nThe apex seal in a Wankel engine is generally a strip of material with a high hardness rating and/or a low friction value when used in the combustion chamber, where it rubs along the interior wall thereof. An example of such a material is graphite, suitably held by a metallic support, whereby both low friction and good material strength are provided, as described in U.S. Pat. No. 3,235,171. Many other materials have been used, among them high-chrome alloys, and the like.\nApex and corner seals serve the function of restricting fluid (gases in the case of rotary (Wankel) engines) to a particular portion of the apparatus. A considerable amount of effort has been expended in the attempt to provide such seals which will maintain a reasonably uniform and effective barrier to the working fluid, typically a gaseous mixture. One approach has been such as is described in U.S. Pat. No. 3,711,229 to Kurio, where the apex seal means is assisted by corner seals which are biased against the chamber wall by springs bearing on the rotary member.\nAnother problem common to seals for rotary engines, pumps, and the like, is the wear on the seal, which allows fluid or gas blow-by and lubricant leakage when it becomes extreme. This in itself would not be a severe drawback, assuming a reasonable life with respect to the running time of the engine or pump, since the problem of seal replacement is of the same approximate magnitude as that of ring replacement in a piston engine or compressor. However, one aspect of seal wear has been manifested in nonuniformity; that is, seals may wear in a fashion which permits the face or edge of the seal to conform to a shape different from the optimum configuration, which, in turn, may lead to pressure loss, excessive oil consumption, or both.\nThe problem of lubrication of seals of the type under discussion is an acute one, presenting considerations which must be taken into account when choosing a seal. Without any lubrication, of course, heat and friction would quickly destroy the seal and severely wear the chamber wall. Lubrication must be supplied, but it must be done in an economical manner, which is to say that a film of oil has to be present, but that a bath of oil would be wasteful.\nSystems are known in this field which compensate for mechanical wear, such as by the use of wedge- or trapezoidal-shaped members, and the like, which displace longitudinally of the seal. There are also means for permitting a certain latitude in manufacturing tolerances. However, there have been no seal systems heretofore capable of compensating well for thermal expansion and contraction with their resulting unpredictable side effects. Necessarily, then, no seal system has yet been able to compensate for all of these variables.\nWith reference particularly to rotary-combustion engines, e.g., the Wankel, it is well-known among those skilled in the art that such engines can produce unacceptably high quantities of noxious emissions under various conditions which can occur readily due either to the construction or the operation of these engines. For example, if an apex seal is used which has gaps in the seal line, blow-by of unburned hydrocarbon results, and these materials then are emitted into the atmosphere unless some sort of post-combustion treatment is applied. Similarly, where gaps occur in the seal line or where the seal is worn from use, the leakage of hydrocarbon reduces combustion efficiency and temperature.\nA further source of hydrocarbon emissions in some rotary-combustion engines is the use of oil-metering arrangements to provide lubrication to apex and corner seals. Because of the need to maintain a film of oil between the chamber wall and the seal, the oil must be supplied in appropriate fashion. In earlier designs, this was done by providing oil under pressure to the seal area itself; however, in most seal configurations, either geometry or economics dictate that the oil be supplied by admixture with the fuel. In these cases, some of the oil tends to be carried out unburned with the exhaust gases, obviously resulting in a detriment to the environment.\nWith all of the problems mentioned hereinabove, a concomitant effect is the lowering of the fuel economy of rotary-combustion engines below the optimum level otherwise possible. With improved seals, blow-by would be decreased, and over-all efficiency increased, leading then to improved fuel economy."}
{"text": "The invention relates to a discharging device for dispensing liquids, more particularly for dispensing pharmaceutical liquids, which discharging device includes a housing, control electronics, a liquid storage receptacle for storing the liquid prior to discharge thereof, a discharge orifice, a conveying device which can be operated manually by means of a first actuator and by means of which the liquid can be conveyed from the liquid storage receptacle to the discharge orifice for the purpose of dispensing the same, and a blocking device that can assume a blocked state, in which operation of the conveying device is mechanically blocked, and a released state, in which operation of the conveying device is possible, for the purpose of preventing another discharging operation from taking place during a blocking period following a discharging operation. For this purpose, the blocking device of a generic discharging device includes an interlock device that is capable of being moved relative to the housing between a first and a second position, the blocking device being configured such that a movement of the interlock device from its first position to its second position causes the blocking device to assume the released state. The interlock device is spring biased by a spring device in the direction of its second position. Furthermore, there is provided a first retaining device, by means of which the interlock device can be retained against the force of this spring device in its first position and which can be released by an electrical signal generated by the control electronics.\nA generic discharging device is disclosed in DE 10 2008 064 559 A1. Provision is made in this device for the mechanical work involved in moving the blocking device between the blocked state and the released state to be imparted to the system by the user during the actuation of the device for the purpose of carrying out a discharging operation. Following the discharging operation, this mechanical work is stored in one or more spring devices held in the tensioned state by a retaining device. This retaining device is formed by a permanent magnet and an electromagnet in the embodiments disclosed in DE 10 2008 064 559 A1, the permanent magnet performing the actual retaining function, whilst it is rendered temporarily ineffective by means of the electromagnet for the purpose of transferring the blocking device to the released state, which is thus achieved. Thus the retaining device is released by activation of the electromagnet.\nThe interlock device of a generic discharging device need not directly serve for blocking the conveying device. Thus both the push rod 78 and the bolt 142 described in DE 10 2008 064 559 A1 are interlock devices in terms of the present invention, since the displacement thereof directly or indirectly results in achieving the released state of the blocking device.\nThe design of the retaining device as a permanent magnet and an electromagnet constitutes a very advantageous solution, but it also involves the problem that strong impacts, caused accidentally or willfully, may cause unwanted movement of the blocking device so as to achieve the released state before the blocking period has elapsed. This poses the risk of administration of an overdose of the active ingredient stored in the liquid storage receptacle."}
{"text": "This invention relates to marine seismic surveying.\nIn order to perform a 3D marine seismic survey, a plurality of seismic streamers, each typically several thousand meters long and containing arrays of hydrophones and associated electronic equipment distributed along its length, are towed at about 5 knots behind a seismic survey vessel, which also tows one or more seismic sources, typically air guns. Acoustic signals produced by the seismic sources are directed down through the water into the earth beneath, where they are reflected from the various strata. The reflected signals are received by the hydrophones in the streamers, digitised and then transmitted to the seismic survey vessel, where they are recorded and at least partially processed with the ultimate aim of building up a representation of the earth strata in the area being surveyed.\nThe Applicant\"\"s current method of towing arrays of such streamers involves towing each streamer by means of its lead-in, ie by means of the armoured electrical or electro-optical cable that supplies control signals and electrical power to, and receives the aforementioned digitised signals from, the streamer, as described in U.S. Pat. No. 4,798,156, the armoured outer sheath of the lead-in, typically made of steel wires or high strength synthetic fibers such as Kevlar, serving to bear the towing forces. Using this method, the Applicant can typically tow a 700 meter wide array of eight streamers, each 4,000 meters long.\nThe drag produced by such an array at the typical towing speed of 5 knots is about 40-45 tonnes, a high proportion of which is cross-line drag due to the transversely extending lead-ins rather than the in-line drag of the streamers themselves. This drag is a very significant contributory factor to the operating costs, primarily fuel costs, of the towing vessel, and therefore to the overall cost of the survey.\nIn order to increase the efficiency and reduce the cost of marine seismic surveys, it is desirable to use even wider arrays containing even more streamers. However, the Applicant has calculated that using its current towing technique, a 1440 meter wide array of ten streamers would produce a drag of over 70 tonnes, which is excessive, and therefore uneconomic.\nIt is an object of the present invention in its principal aspects to alleviate this problem.\nThe lead-ins, which because of their combined signal and power transmission and towing functions are rather expensive, are sometimes damaged in use, eg by shark bites or by impacts with debris. Currently, it is difficult if not almost impossible to repair a damaged lead-in, especially while at sea. It is an object of the present invention in its final aspect to provide solution to this problem.\nAccording to a first aspect of the present invention, there is provided a method of performing a marine seismic survey with an array of seismic streamers towed behind a survey vessel, wherein two or more of the streamers are towed from a single lead-in which comprises a load-bearing outer sheath to bear the towing forces, an electrical or electro-optical core via which control signals and electrical power are supplied to, and data signals are received from, the streamers being towed by the lead-in, and a connector device series connected in the lead-in, the connector device comprising a body member mechanically coupled between the respective portions of the load-bearing outer sheath of the lead-in on either side thereof, to transmit towing forces therebetween, and a towing attachment for mechanical connection to one of the two or more streamers.\nAccording to a second aspect of the invention, there is provided an array of seismic streamers adapted to be towed behind a survey vessel, the array comprising at least one lead-in which has two or more of the streamers connected to be towed thereby, and which comprises a load-bearing outer sheath to bear the towing forces, an electrical or electro-optical core via which control signals and electrical power are supplied to, and data signals are received from, the streamers being towed by the lead-in, and a connector device series connected in the lead-in, the connector device comprising a body member mechanically coupled between the respective portions of the load-bearing outer sheath of the lead-in on either side thereof, to transmit towing forces therebetween, and a towing attachment for mechanical connection to one of the two or more streamers.\nIn a preferred implementation of the invention, the array comprises 2N streamers, where N is at least 2, the streamers are substantially uniformly spaced apart and symmetrically distributed on each side of the centre line of the vessel, and at least two streamers on each side of the centre line of the vessel are towed from a respective such lead-in.\nThe Applicant has found that by use of selected implementations of the method of the invention, the drag produced by the aforementioned 1440 meter wide ten streamer array can be reduced from over 70 tonnes to a level in the region of 50 tonnes.\nThe load-bearing outer sheath of the or each such lead-in is preferably made from high strength synthetic fibers, such as Kevlar fibres.\nPreferably, the inner core of the lead-in at each end of the or each connector device is terminated in a respective watertight termination, and the connector device includes elongate flexible electrical or electro-optical connector means for establishing electrical or electro-optical connection between said watertight terminations and between one of said terminations and the streamer connected to the towing attachment.\nConveniently, the inner core of the part of the lead-in on the vessel side of each connector includes separate electrical conductors or optical fibres for the control and data signals of each streamer connected directly or indirectly to that part via that connector.\nAccording to a third aspect of the present invention, there is provided a connector device adapted for series connection in a lead-in to permit the lead-in to tow at least two seismic streamers, the lead-in comprising a load-bearing outer sheath, to bear the towing forces, and an electrical or electro-optical core via which control signals and electrical power are supplied to, and data signals are received from, the streamers, the connector device comprising a body member adapted to be mechanically coupled between the respective portions of the load-bearing outer sheath of the lead-in on either side of the body member, to transmit towing forces therebetween, and a towing attachment for mechanical connection to one of the two or more streamers, said towing attachment being secured to the body member.\nThe or each body member may advantageously comprise at least two generally tubular portions interconnected by a universal joint, the towing attachment being secured to one of the tubular body portions.\nIn a preferred implementation of the invention, the or each body member comprises three tubular body portions interconnected by two universal joints, the towing attachment being secured to the middle one of tubular body portions.\nThe or each towing attachment preferably includes clamping means for detachably clamping it to the tubular body portion to which it is secured.\nAdvantageously, the or each towing attachment ha s a towing bracket pivotally connected to it, said towing bracket being pivotal about an axis generally perpendicular to the axis of the body member.\nAccording to a final aspect of the present invention, there is provided a connector device adapted for series connection in a lead-in used for towing a streamer, the lead-in comprising a load bearing outer sheath to bear the towing forces, and an electrical or electro-optical core via which control signals and electrical power are supplied to, and data signals are received from, the streamer, the connector device comprising a tubular body member which is adapted to be mechanically coupled between the respective portions of the load-bearing outer sheath of the lead-in on either side of the body member, to transmit towing forces therebetween, and which preferably comprises at least two generally tubular portions interconnected by a universal joint.\nPreferably, the or each universal joint of the or each body member is a Cardan-type universal joint.\nIn all aspects of the invention, at least part of the or each body member is preferably made from titanium.\nThe invention will now be described, by way of example only, with reference to the accompanying drawings, of which:"}
{"text": "1. Field of the Invention\nThis invention relates to apparatus for supporting a conduit from a beam and more particularly to a beam clamp coated with a moisture, resistant polymeric material in which a clamping member is secured to a U-bolt that supports a conduit adjacent to the beam with polymeric encapsulated nuts engaging the threads of the U-bolt is sealed relation to prevent corrosion of the threaded connection.\n2. Description of the Prior Art\nBeam clamps are well known in the art of securing a conduit, pipe, or the like in a preselected position to a structural member, such as an I-beam. A variety of types of beam clamps are known to position the conduit either vertically, parallel, or on edge relative to the longitudinally extending horizontal flange of the I-beam. U.S. Pat. No. 2,338,006 is an example of a clamp device releasably engageable with the vertical edge of a supporting structure to position the conduit in spaced parallel relation to the supporting structure. This type of clamp employs a threaded member for securing a U-shaped strap by a clip to the wall of the supporting structure. The clip is secured to the supporting structure by the threaded connection of a bolt and nut. The nut engages the threaded portion of the bolt that extends through the supporting structure.\nOne of the disadvantages of this type of arrangement is that threaded connections are subject to corrosive damage when employed in a corrosive environment. In many cases a beam clamp is located in a corrosive environment where moisture can easily enter the threaded connection of nuts and bolts. The corrosive damage can \"freeze\" the threaded connection. Consequently the beam clamp can become rendered useless by creating diffucult maintenance problems in disengaging the nut from the bolt to disassemble the clamp from the beam and the conduit.\nIt is also well known in the art as disclosed in U.S. Pat. Nos. 3,724,706; 3,784,236; and 3,799,584 to coat components of an assembly that must be capable of efficient assembly and disassembly in a corrosive environment with a moisture resistant, insulating, resilient, polymeric material, such as a polyvinylchloride (PVC). By providing male and female coupling members with a moisture, resistant polymeric coating a moisture resistant seal is formed around the otherwise exposed threads of the coupling members. In this manner moisture is prevented from entering the point of engagement, for example, of a nut on a threaded bolt to prevent freezing of the nut on the bolt.\nWhile it is known to coat the components of a beam clamp, i.e. the clamping member and the U-bolt, with a moisture resistant, polymeric coating, the threaded connection of the clamping member to the U-bolt is exposed. The components generally are assembled in the field where the nuts and bolts are exposed to the affects of moisture.\nTo resist the deleterious affects of moisture it is known to apply a moisture resistant, polymeric material by brush coating the exposed nuts and the U-bolt. However, a moisture resistant coating applied in the manner to the exposed threaded connections and fasteners in a corrosive environment has proved unsatisfactory in preventing corrosion. Brushing on a polymeric coating after installation will not deter the accummulated affects of the corrosive atmosphere prior to the application of the coating. Furthermore, if improper attention is given to the brushing application of polymeric material to the exposed threads and nuts, an insufficient coating thickness to prevent corrosion may result.\nTherefore, there is need to provide for a beam clamp a connection of the clamping member and the conduit support member capable of preventing corrosive damage while permitting rapid assembly and disassembly of the beam clamp components."}
{"text": "1. Field of the Invention\nThe present invention relates generally to non-volatile magnetic memory and particular to multi-state magnetic memory.\n2. Description of the Prior Art\nComputers conventionally use rotating magnetic media, such as hard disk drives (HDDs), for data storage. Though widely used and commonly accepted, such media suffer from a variety of deficiencies, such as access latency, higher power dissipation, large physical size and inability to withstand any physical shock. Thus, there is a need for a new type of storage device devoid of such drawbacks.\nOther dominant storage devices are dynamic random access memory (DRAM) and static RAM (SRAM), which are volatile and very costly but have fast random read/write access time. Solid state storage, such as solid-state-nonvolatile-memory (SSNVM) devices having memory structures made of NOR/NAND-based Flash memory, providing fast access time, increased input/output (TOP) speed, decreased power dissipation and physical size and increased reliability but at a higher cost which tends to be generally multiple times higher than hard disk drives (HDDs).\nAlthough NAND-based Flash memory is more costly than HDD's, it has replaced magnetic hard drives in many applications such as digital cameras, MP3-players, cell phones, and hand held multimedia devices due, at least in part, to its characteristic of being able to retain data even when power is disconnected. However, as memory dimension requirements are dictating decreased sizes, scalability is becoming an issue because the designs of NAND-based Flash memory and DRAM memory are becoming difficult to scale with smaller dimensions. For example, NAND-based Flash memory has issues related to capacitive coupling, few electrons/bit, poor error-rate performance and reduced reliability due to decreased read-write endurance. Read-write endurance refers to the number of reading, writing and erase cycles before the memory starts to degrade in performance due primarily to the high voltages required in the program, erase cycles.\nIt is believed that NAND Flash, especially multi-bit designs thereof, would be extremely difficult to scale below 45 nanometers. Likewise, DRAM has issues related to scaling of the trench capacitors leading to very complex designs that are becoming increasingly difficult to manufacture, leading to higher cost.\nCurrently, applications commonly employ combinations of EEPROM/NOR, NAND, HDD, and DRAM as a part of the memory in a system design. Design of different memory technology in a product adds to design complexity, time to market and increased costs. For example, in hand-held multi-media applications incorporating various memory technologies, such as NAND Flash, DRAM and EEPROM/NOR Flash memory, complexity of design is increased as are manufacturing costs and time to market. Another disadvantage is the increase in size of a device that incorporates all of these types of memories therein.\nThere has been an extensive effort in development of alternative technologies such as Ovanic RAM (or phase-change memory), Ferromagnetic RAM (FeRAM), Magnetic RAM (MRAM), probe-based storage such as Millipede from International Business Machines, Inc. of San Jose, Calif., or Nanochip, and others to replace memories used in current designs such as DRAM, SRAM, EEPROM/NOR Flash, NAND Flash and HDD in one form or another. Although these various memory/storage technologies have created many challenges, there have been advances made in this field in recent years. MRAM seems to lead the way in terms of its progress in the past few years to replace all types of memories in the system as a universal memory solution.\nOne of the problems with prior art memory structures is that the current and power requirements are too high to make a functional memory device or cell. This also poses a key concern regarding the reliability of such devices due to likely dielectric breakdown of the tunneling barrier layer and thereby making it non-functional.\nThe challenge with other prior art techniques has been that the switching current is too high to allow the making of a functional device for memory applications due to the memory's high power consumption. Several recent publications, such as those cited below as references 5 and 6 (5′6) have shown that the switching current can be reduced by having the memory element pinned by two anti-ferromagnetic (AF)-couple layers resulting in spin oscillations or “pumping” and thereby reducing the switching current.\nAn additionally known problem is using magnetic memory to store more than two states therein. To this end, multi-level or multi-state magnetic memory cells or elements for storing more than one bit of information do not exist.\nWhat is needed is magnetic memory for storing more than one bit of digital information."}
{"text": "The present invention relates to a method of cooling retort or chamber coke.\nAside from the conventional wet quenching where, as a rule, the sensible heat of the coke is not utilized at all, or recovered only to a very small extent, dry cooling of coke is known in which the heat of the retort or chamber coke having a temperature of about 1,000.degree. C. is removed by a circulated gas and utilized in a waste heat boiler for producing steam. It is further known to use this steam for preheating the moist coking coal, with the heat being transferred to the coal, for example, in a fluidized bed or rotary drum drier through suitable heating surfaces (see German patent application Nos. P 30 13 325.6 and P 30 34 952, U.S. Pat. Nos. 4,354,903 and 4,422,858, respectively, which are here incorporated by reference).\nIt is also known to supply the chamber coke having a temperature of about 1,000.degree. C. to a first cooling and gasification stage for partial degasification under the addition of steam and to cool it down to 500.degree.-800.degree. C., and then to direct it through a pressure metering system to a pressure-type gasifier where it is gasified by fixedbed gasification while adding air or oxygen, and steam (German application No. P 30 32 212.4 or U.S. Pat. No. 4,422,858 also incorporated by reference).\nWith such coke cooling methods, the produced steam can only be employed for preheating the coal to be coked, in view of the heat economy of the coking plant. Exceptionally, a power plant may be in the neighborhood of the coking plant, where the steam from the dry coke cooling may be used for generating power. However, this is not the case as a rule. Further, if the steam is employed for preheating moist coking coal, as usual, only superheated steam is used for transferring the heat to the coal through indirectly heated surfaces. This means that, as compared to the relatively high temperature of about 1,000.degree. C. of the coke coming from the chambers, the heat is utilized at a relatively low temperature level (t=232.degree. C. at 30 bar, or t=249.degree. C. at 40 bar) large heat exchanger surfaces or, in addition, special and very expensive measures must be provided (such as a fluidized bed) for improving the heat transfer."}
{"text": "This invention relates specifically to a wheeled apparatus for transporting skis and related equipment.\nSkiers need to carry lots of equipment when they go to the slopes. If these people are parents with smaller children then they must also carry the little one's share. This can take lots of the fun out of the day's adventure. It can also be dangerous what with whirling skis and sharp ski poles. It seems that there is a need for a device which can safely transport a variety of ski equipment and be conveniently stored away in a very small area.\nParker, U.S. Pat. No. 5,340,153, discloses a ski carrier that will carry two pair of skis on a wheeled dolly similar to the type used to transport utility poles behind a truck. That is, attach one end of the long object to a pair of wheels and attach the opposite end to a vehicle.\nGarvey Jr, U.S. Pat. Nos. 4,666,184 and 4,792,159 discloses a two wheeled carrier for a single pair of skis and Kybutz, U.S. Pat. No. 4,540,198 discloses a one wheel device for a single pair of skis. All of these devices are adequate for their intended uses but fall far short of the goal of this invention; to provide a multi-purpose cart that can not only carry three or more pair of skis, boots and poles but can also be used to carry other gear."}
{"text": "This invention relates to methods and apparatus for parsing Extensible Markup Language (XML) data streams.\nExtensible Markup Language (XML) is a meta-markup language that provides a format for describing structured data. XML is similar to HTML in that it is a tag-based language. By virtue of its tag-based nature, XML defines a strict tree structure or hierarchy. XML is a derivative of Standard Generalized Markup Language (SGML) that provides a uniform method for describing and exchanging structured data in an open, text-based format. XML utilizes the concepts of elements and namespaces. Compared to HTML, which is a display-oriented markup language, XML is a general-purpose language used to represent structured data without including information that describes how to format the data for display.\nXML xe2x80x9celementsxe2x80x9d are structural constructs that consist of a start tag, an end or close tag, and the information or content that is contained between the tags. A xe2x80x9cstart tagxe2x80x9d is formatted as xe2x80x9c less than tagname greater than xe2x80x9d and an xe2x80x9cend tagxe2x80x9d is formatted as xe2x80x9c less than /tagname greater than xe2x80x9d. In an XML document, start and end tags can be nested within other start and end tags. All elements that occur within a particular element must have their start and end tags occur before the end tag of that particular element. This defines a strict tree-like structure. Each element forms a node in this tree, and potentially has xe2x80x9cchildxe2x80x9d or xe2x80x9cbranchxe2x80x9d nodes. The child nodes represent any XML elements that occur between the start and end tags of the xe2x80x9cparentxe2x80x9d node. XML accommodates an infinite number of database schemas. Within each schema, a xe2x80x9cdictionaryxe2x80x9d of element names is defined. The dictionary of element names defined by a schema is referred to as a xe2x80x9cnamespace.xe2x80x9d Within an XML document, element names are qualified by namespace identifiers. When qualified by a namespace identifer, a tag name appears in the form xe2x80x9c[namespace]:[tagname]xe2x80x9d. This model enables the same element name to appear in multiple schemas, or namespaces, and for instances of these duplicate element names to appear in the same XML document without colliding.\nStart tags can declare an arbitrary number of xe2x80x9cattributesxe2x80x9d which declare xe2x80x9cproperty valuesxe2x80x9d associated with the element being declared. Attributes are declared within the start tag using the form xe2x80x9c less than [tagname] [attributel],[attribute2] . . . ,[attributeN] greater than xe2x80x9d, where an attribute through attributeN are declarations of an arbitrary number of tag attributes. Each attribute declaration is of the form xe2x80x9c[attributeName]=[attributeValue]xe2x80x9d where each attribute is identified by a unique name followed by an xe2x80x9c=xe2x80x9d character, followed by the value of the attribute.\nWithin an XML document, namespace declarations occur as attributes of start tags. Namespace declarations are of the form xe2x80x9cxmlns:[prefix]=[uri]xe2x80x9d. A namespace declaration indicates that the XML document contains element names that are defined within a specified namespace or schema. Prefix is an arbitrary designation that will be used later in the XML document as an indication that an element name is a member of the namespace declared by uri. The prefix is valid only within the context of the specific XML document. xe2x80x9cUrixe2x80x9d or universal resource indicator is either a path to a document describing a specific namespace or schema or a globally unique identifier of a specific namespace or schema. Uri is valid across all XML documents. Namespace declarations are xe2x80x9cinheritedxe2x80x9d, which means that a namespace declaration applies to the element in which it was declared as well as to all elements contained within that element.\nNamespace inheritance within an XML document allows non-qualified names to use xe2x80x9cdefaultxe2x80x9d namespaces. Default namespaces are explicitly declared as attributes of start tags. Default namespace declarations are of the form xe2x80x9cxmlns=[uri]xe2x80x9d. Note that the declaration of a default namespace is equivalent to the declaration of a non-default namespace but the prefix is omitted A namespace specification within an XML document is said to have a xe2x80x9cscopexe2x80x9d which includes all child nodes beneath the namespace specification.\nOne exemplary usage of XML is the exchange of data between different entities, such as client and server computers, in the form of requests and responses. A client might generate a request for information or a request for a certain server action, and a server might generate a response to the client that contains the information or confirms whether the certain action has been performed. The contents of these requests and responses are xe2x80x9cXML documentsxe2x80x9d, which are sequences of characters that comply with the specification of XML. Part of the document exchange process between clients and servers involves parsing the XML documents when they are received. In many cases, it is convenient to represent these XML documents in memory as a hierarchical tree structure. Once the hierarchical tree structure is built, the actual parsing process can begin. Consider the following exemplary XML code:\nThis code includes three XML namespace declarations that are each designated with xe2x80x9cxmlnsxe2x80x9d. The declarations include a prefix, e.g. xe2x80x9cpersonxe2x80x9d, xe2x80x9cdsigxe2x80x9d, and xe2x80x9ctransxe2x80x9d respectively, and the expanded namespace to which each prefix refers, e.g. www.schemas.org/people, dsig.org, and www.schemas.org/transactions respectively. This code tells any reader that if an element name begins with xe2x80x9cdsig:xe2x80x9d its meaning is defined by whoever owns the www.dsig.org namespace. Similarly, elements beginning with the xe2x80x9cperson:xe2x80x9d prefix have meanings defined by the www.schemas.org/people namespace and elements beginning with the xe2x80x9ctransxe2x80x9d prefix have meanings defined by the www.schemas.org/transactions namespace. It is important to note that another XML document that incorporated elements from any of the namespaces included in this sample might declare prefixes that are different from those used in this example. As noted earlier, prefixes are arbitrarily defined by the document author and have meaning only within the context of the specific element of the specific document in which they are declared.\nNamespaces ensure that element names do not conflict, and clarify who defined which term. They do not give instructions on how to process the elements. Readers still need to know what the elements mean and decide how to process them. Namespaces simply keep the names straight.\nFIG. 1 shows how the structure of the above code can be represented in a hierarchical tree structure. In FIG. 1, all of the elements or nodes are set out in an exemplary tree that represents the XML document. Such a structure is typically constructed in memory, with each node containing all data necessary for the start and end tags of that node.\nIt has been typical in the past to build the entire tree structure, such as the one shown in FIG. 1, before parsing the XML document. For large XML documents, this can consume a great deal of memory and processor time. Thus, it would be desirable to avoid this process if at all possible.\nXML parsers are used to various applications to process XML documents. Parsers must know what particular elements mean and how to process them. Tags from multiple namespaces can be mixed, which is essential with data coming from multiple sources across the Web. With namespaces, both elements could exist in the same XML-based document instance but could refer back to two different schemas, uniquely qualifying their semantics. Parsers typically take the form of a code library that can be used by developers in conjunction with higher level languages such as C++ or Java. Using functions provided by such a code library, developers can access the structure of an XML document, enumerate its elements and their attributes, and manipulate the information that is contained within the document\"\"s prolog. A simple example would be an XML parser utility that checks for xe2x80x9cwell-formedxe2x80x9d or xe2x80x9cvalidxe2x80x9d documents, and serves as the equivalent of an syntax checker.\nXML parsers typically read XML files or data streams and construct a hierarchically structured tree, such as the one appearing in FIG. 1, as a data structure in memory. The XML parser then typically hands off this data structure data to viewers and other applications for processing. So, in the example XML code discussed above, a parser would first build the entire tree structure that is shown in FIG. 1 prior to interpreting the contents of the document. Only after the entire tree structure was built in memory would the parser begin to interpret the document.\nOne problem that is associated with XML parsers such as this is that they have to build an entire hierarchically structured tree in memory before interpreting the contents of the document. This approach is not efficient because of the demands it places on the memory that is required to store the tree structure and the speed with which information can be conveyed to a client. For example, this type of approach is not efficient for an application that is doing work in connection with a large quantity of XML data that might be streaming in at a relatively slow speed. Consider, for example, that a client asks a server for a list of all messages of a certain type that are in a certain folder. The entire message list is going to be returned by the server as one large data stream. If the client has to wait for the entire message list to be returned from the server, then the client cannot begin to display any portion of the list until all of the data has been received. Furthermore, this process requires the parser to make at least two passes over the data; the first pass to build the tree structure, and the second pass to traverse the nodes of the tree to interpret the contents of the document. This approach requires a large memory overhead (for storing the XML data and building the hierarchical tree structure) which, in turn, impacts the speed with which responses can be used by client applications.\nThis invention arose out of concerns associated with providing improved XML parsers and methods of parsing XML data streams that reduce memory overhead and increase the speed with which XML data can be provided and used by a client\nVarious features of the invention provide methods and systems for parsing an XML data stream that do not require that an entire hierarchical tree structure be built and stored in memory in conjunction with parsing activities.\nVarious parsing support structures are defined that enable the inventive parsing activities. An element stack is defined and used as a way to organize parsing activities and maintain a definable place within the structure of an XML data stream as it is received. During parsing activities the element stack grows and shrinks in connection with how deep into an XML data stream the parser goes.\nA namespace hierarchy is defined as the XML data stream is received. In one described embodiment, the namespace hierarchy is a collection of namespace objects that each represent a namespace specification that is encountered in the XML data stream. Each object includes a namespace specification and, in the case of non-default namespace, a namespace prefix. Within the hierarchy, namespace objects refer to their parent. The namespace hierarchy is used for mapping a particular encountered namespace specification into a unique value that represents both the namespace specification and an element tag in which the namespace occurs.\nA dictionary collection is defined and includes one or more dictionaries. Each dictionary is specifically associated with a namespace specification that is encountered in the XML data stream. The dictionaries contain one or more entries for one or more tag names, and a globally unique token that is associated with each tag name. The token is returned and placed on the element stack, along with another special value that enables the proper state to be maintained during processing of the XML data stream.\nWhen an XML element is encountered, the element stack is used to organize parsing activities as the parser makes its way through the XML data stream. During parsing activities, various notifications are generated by the parser and sent to the client or application when various events occur. Exemplary notifications include, without limitation, those that are sent when: a new tag is encountered, the children of a tag are entered, data is encountered within a tag, and a close tag is encountered. These notifications can be used to organize parsing activities. For example, when notification is received that a new tag has been encountered, the element can be pushed onto the element stack. Similarly, when notification is received that a close tag has been encountered, the top frame can be removed from the element stack."}
{"text": "There are many reasons for wanting to incorporate antimicrobial agents into textiles. For instance, antimicrobial activity may enhance the durability of textiles by inhibiting the growth of bacteria and fungi that live in textiles and cause damage to their components. Textiles with antimicrobial activity also prevent the creation of odors by preventing or reducing microbes from feeding on the organic materials in the textile. Additionally, antimicrobial treated textiles may kill pathogenic microbes thereby protecting the wearer from exposure to disease agents.\nOther uses of antimicrobial treated textiles include wound care. For example, a medical bandage that has been treated with antimicrobial agents offers the wound a favorable environment for healing and further prevents bacterial organisms from growing at the site. Furthermore, antimicrobial cleaning cloths may prevent disease-causing microorganisms from surviving or even proliferating in the internal crevices of the textile when the cloth soaks up contaminated fluid.\nAntimicrobial bedding linens offer an enhanced level of hygiene and security in hospital or hotel rooms. Their use may limit the transmission of bacterial disease in hospitals and lower the rates at which hospitalized patients acquire nosocomial infections, which have been increasing in frequency in recent years.\nWhile textiles possessing antimicrobial characteristics have been introduced in recent years, most of these textiles are prepared by adding antimicrobial agents to the textile. Specifically, antimicrobial agents are impregnated into or coated onto the fibers of the textile. The antimicrobial agents are not, however, permanently bound to the textile fibers and are thus prone to leaching, such that these agents are absorbed by human skin. In addition to leaching, the antimicrobial agents are not able to withstand numerous washings, so that the antimicrobial effect may not last long. While certain N-halamine treated textiles are known to be regenerable after exhaustion by the treated textiles, use of chlorine bleach is necessary to regenerate antimicrobial activity.\nThus, there is a need to overcome the aforementioned disadvantages. Current antimicrobial treatment of textile leaves a need for a means of providing various textiles with minimal leaching, more permanent antimicrobial characteristics, and that does not require use of chlorine bleach for regenerable antimicrobial properties."}
{"text": "Color-forming coupler compounds which react during photographic development with the development product of aromatic amino developing agents to form color images are well known. Generally, these color-forming compounds are colorless or substantially colorless. This lack of color is usually desirable when the coupler is to be incorporated in the emulsion layer and the unused coupler remains after formation of the colored image.\nUNFORTUNATELY, THE ABSORPTION SPECTRA OF CONVENTIONAL DYES FORMED FROM COLOR-FORMING COUPLER MATERIALS ARE NEVER \"CLEAN\". Thus, invariably, cyan dyes exhibit various degrees of unwanted absorption in the green and blue regions of the spectrum, while magenta dyes exhibit unwanted absorption in the blue region and to some degree also in the red region of the spectrum. The term \"unwanted absorption\" is intended herein to mean, with respect to a cyan dye, measurable absorption in the green and blue regions of the visible spectrum. Unwanted absorption is illustrated in FIG. 1, as that portion of the curve above the base line at wavelengths lower than about 600 millimicrons.\nWhereas the human eye and brain have the capability of automatically compensating for a certain amount of such unwanted absorption in positive color photographs and transparencies which contain such conventional dyes, machines and equipment for reproducing photographs and transparencies do not have such compensating capability. Hence, photographic materials which must be mechanically reproduced, for example, from color negatives or positive transparencies, must be color-corrected so as to largely overcome the effects of unwanted absorption. Color correction becomes increasingly necessary when it is desired to reproduce the color element by means of a series of negative-positive-negative-positive steps such as those usually used in the preparation of (from valuable \"master\" films) motion picture films for release.\nOne very useful method of compensating for the unwanted absorption caused by a dye derived from a given dye-forming coupler involves the inclusion of a special colored coupler into the photographic element (in addition to the image-forming coupler that is to be \"color-corrected\"). The colored coupler absorbs both green and blue light, in the case of correcting for unwanted absorption in a cyan dye. This colored coupler theoretically is capable of reacting with oxidized color developer (during the color development processing step) to yield the usable cyan image dye while simultaneously losing its ability, in proportion to development, to absorb in the green and blue regions of the spectrum, thereby \"correcting\" for the unwanted green and blue absorption of the dye derived from the major cyan dye-forming coupler in the photographic element. (See Chapter 13, \"Masking and Coloured Couplers\" of R. W. G. Hunt's book, The Reproduction of Colour, pages 233, 263, published by John Wiley & Sons, 1967).\nCouplers which are in themselves more or less strongly colored by virtue of containing a chromophore group which is split off or destroyed during and by means of the coupling reaction with the result that the original color of the coupler is destroyed and a new dye is formed upon coupling are disclosed in U.S. Pat. Nos. 2,453,661, Nov. 9, 1948; 2,445,169, Nov. 30, 1948; 2,455,170, Nov. 30, 1948; 2,521,908, Sept. 12, 1950; and 2,706,684, Apr. 19, 1955.\nUpon development of an emulsion layer containing one of these colored couplers, the original color of the colored coupler is destroyed and a new color is formed by the coupling reaction at those points where development occurs. There is formed in this way a composite dye image consisting of the new dye and the residual colored coupler. The theory of color correction requires that the sum of the absorption of the residual colored coupler and the undesired absorption of the image dye should be as constant as possible. Conventional color correction to date has resulted either in non-ideal (non-uniform) correction, or uniform correction with somewhat higher Dmin than is desirable. (The higher Dmin apparently results from the necessity to use relatively large amounts of certain color-correcting couplers which have relatively low activity.) Another shortcoming that exists with respect to certain of the conventional color-correcting couplers that otherwise exhibit good correction relates to their color prior to their reaction with oxidized color developer. The peak absorption of such colored couplers was often too hypsochromic, by several wavelengths, and thus yielded corrections which were somewhat less than desirable.\nThe use of certain prior art color-correcting couplers, such as those described in the Examples, below, yield excellent dyes, but are not capable of yielding \"ideal\" or \"uniform\" color-correction,. This fact is illustrated by FIGS. 3 and 4. Note that substantially uniform absorption in the green and blue regions (particularly in the green) is not achieved by the materials used in the preparation of FIGS. 3 and 4. Ideal color correction should result in perfectly level \"G\" and \"B\" curves in the Figures, such as those marked \"ideal\" in FIGS. 3 and 4.\nUniform color-correction throughout the development of the color emulsion has been extremely difficult to achieve because of the apparent requirement of uniformity or similarity in each of the many different chemical and physical parameters involved in the color-correcting process. For example, the reaction rates and relative reactivities of each of the couplers in any \"color correction\" combination must be matched at every stage of the development step, even though the concentration of developer, pH, concentration of reaction by-products and many other aspects of the chemical color development process not only differ widely from one photo-finisher to another, but also differ widely during any single color development operation."}
{"text": "Wireless communication systems are widely deployed to provide various types of communication content such as voice, data, and so on. These systems may be multiple-access systems capable of supporting communication with multiple users by sharing the available system resources (e.g., bandwidth and transmit power). Examples of such multiple access systems include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, 3GPP Long Term Evolution (LTE) systems, and orthogonal frequency division multiple access (OFDMA) systems.\nGenerally, a wireless multiple-access communication system can simultaneously support communication for multiple wireless terminals. Each terminal communicates with one or more base stations via transmissions on the forward and reverse links. The forward link (or downlink) refers to the communication link from the base stations to the terminals, and the reverse link (or uplink) refers to the communication link from the terminals to the base stations. This communication link may be established via a single-in-single-out, multiple-in-single-out or a multiple-in-multiple out (MIMO) system.\nSome conventional advanced devices include multiple radios for transmitting/receiving using different Radio Access Technologies (RATs). Examples of RATs include, e.g., Universal Mobile Telecommunications System (UMTS), Global System for Mobile Communications (GSM), cdma2000, WiMAX, WLAN (e.g., WiFi), Bluetooth, LTE, and the like.\nAn example mobile device includes an LTE User Equipment (UE), such as a fourth generation (4G) mobile phone. Such 4G phone may include various radios to provide a variety of functions for the user. For purposes of this example, the 4G phone includes an LTE radio for voice and data, an IEEE 802.11 (WiFi) radio, a position location e.g., Global Positioning System (GPS)) radio, and a Bluetooth radio, where two of the above or all four may operate simultaneously. While the different radios provide useful functionalities for the phone, their inclusion in a single device gives rise to coexistence issues. Specifically, operation of one radio may in some cases interfere with operation of another radio through radiative, conductive, and/or resource collision, other interference mechanisms. Coexistence issues include such interference.\nThis is especially true for the LTE uplink channel, which is adjacent to the Industrial Scientific and Medical (ISM) band and may cause interference therewith It is noted that Bluetooth and some Wireless LAN (WLAN) channels fall within the ISM band. In some instances, a Bluetooth error rate can become unacceptable when LTE is active in some channels of Band 7 or even Band 40 for some Bluetooth channel conditions. Even though there is no significant degradation to LTE, simultaneous operation with Bluetooth can result in disruption in voice services terminating in a Bluetooth headset. Such disruption may be unacceptable to the consumer. A similar issue exists when LTE transmissions interfere with position location. Currently, there is no mechanism that can solve this issue because LTE by itself does not experience any degradation.\nWith reference specifically to LTE, it is noted that a UE communicates with an evolved NodeB (eNB; e.g., a base station for a wireless communications network) to inform the eNB of interference seen by the UE on the downlink. Furthermore, the eNB may be able to estimate interference at the UE using a downlink error rate. In some instances, the eNB and the UE can cooperate to find a solution that reduces interference at the UE, even interference due to radios within the UE itself. However, in conventional LTE, the interference estimates regarding the downlink may not be adequate to comprehensively address interference.\nIn one instance, an LTE uplink signal interferes with a Bluetooth signal or WLAN signal. However, such interference is not reflected in the downlink measurement reports at the eNB. As a result, unilateral action on the part of the UE (e.g., moving the uplink signal to a different channel) may be thwarted by the eNB, which is not aware of the uplink coexistence issue and seeks to undo the unilateral action. For instance, even if the UE re-establishes the connection on a different frequency channel, the network can still handover the UE back to the original frequency channel that was corrupted by the in-device interference. This is a likely scenario because the desired signal strength on the corrupted channel may sometimes be higher than reflected in the measurement reports of the new channel based on Reference Signal Received Power (RSRP) to the eNB. Hence, a ping-pong effect of being transferred back and forth between the corrupted channel and the desired channel can happen if the eNB uses RSRP reports to inform handover decisions.\nOther unilateral action on the part of the UE, such as simply stopping uplink communications without coordination of the eNB may cause power loop malfunctions at the eNB. Additional issues that exist in conventional LTE include a general lack of ability on the part of the UE to suggest desired configurations as an alternative to configurations that have coexistence issues. For at least these reasons, uplink coexistence issues at the UE may remain unresolved for a long time period, degrading performance and efficiency of other radios of the UE."}
{"text": "In the prior art, automatic transmission, particularly for motor vehicles, comprise planetary gear sets shifted by means of friction and shifting elements such as clutches and brakes and usually have one starting element, such as a hydrodynamic converter or a fluid clutch, optionally provided with a lock-up clutch and subject to a slip effect.\nSuch a transmission results from EP 0 434 525 A1. It essentially comprises one drive shaft and one driven shaft disposed parallel with each other, one double planetary gear set disposed concentrically with the driven shaft and five shifting elements in the form of three clutches and two brakes, the optional locking of which by pairs determines the different gear ratios between the drive shaft and the driven shaft. This transmission has a front-mounted range change set and two power paths so that six forward gears are obtained by selective engagement by pairs of the five shifting elements.\nHere two clutches are needed in the first power path for transmitting the torque from the front-mounted range change set to two elements of the double planetary gear set. These are situated in the power flow direction essential downstream of the front-mounted range change set in the direction of the double planetary gear set. In the second power path, one other clutch is provided which loosely connects it with another element of the double planetary gear set. The clutches are arranged here so that the internal disc carrier constitutes the output.\nIn addition, the U.S. Pat. No. 6,139,463 publication has disclosed a compact multiple stage transmission of a planetary design, particularly for motor vehicles, which has two planetary gear sets and one front-mounted range change set, the same as three clutches and two brakes. In this known multiple stage transmission, two clutches C-1 and C-3 are provided in a first power path for transmitting the torque from the front-mounted range change set to the two planetary gear sets. The external disc carrier or the cylinder or piston and pressure compensation side of the clutch C-3 is here connected with a first brake B-1. Besides, the internal disc carrier of the third clutch C-3 is connected with the cylinder or piston and pressure compensation side of the first clutch C-1, the internal disc carrier of the first clutch C-1 is located on the output side and connected with a sun gear of the third planetary gear set.\nFrom the Applicant's DE 199 49 507 A1 is further known a multiple stage transmission where two non-shiftable front-mounted range change sets are provided on the drive shaft which, on the output side, produce two rotational speeds which, together with the rotational speed of the drive shaft, are optionally shiftable by selective engaging of the used shifting elements, to a shiftable double planetary gear set acting upon the driven shaft in a manner such that to change over from one gear to the next sequential higher or lower gear of the two specifically actuated shifting elements, respectively, only one shifting element has to be engaged or disengaged.\nFrom DE 199 12 480 A1 is known an automatically shiftable motor vehicle transmission having three spider planetary gear sets, the same as three brakes and two clutches, for shifting six forward gears and one reverse gear and having one drive shaft and one driven shaft. The automatically shiftable motor vehicle transmission is designed so that the drive shaft is directly connected with the sun gear of the second planetary gear set and the drive shaft is connectable, via the first clutch, with the sun gear of the first planetary gear set and/or via the second clutch with the spider of the first planetary gear set. Additionally or alternatively, the sun gear of the first planetary gear set is connectable, via the first brake, with the housing of the transmission and/or the spider of the first planetary gear set with the housing, via the second brake, and/or the sun gear of the third planetary gear set with the housing via the third brake.\nDE 102 13 820 A1 has also disclosed a multiple stage automatic transmission comprising one first input path T1 of a first reduction ratio; one input path T2 which has a higher reduction ratio than said input path T1; one planetary gear set having four elements, said four elements being one first element, one second element, one third element and one fourth element in the sequence of elements in a rotational speed diagram; one clutch C-2 which transmits rotation of the input path T2 to the first element S3; and clutch C-1 which transmits the rotation from the input path T2 to the fourth element S2; one clutch C-4 which transmits rotation from the input path T1 to the first element; one clutch C-3 which transmits rotation from the input path T1 to the second element C-3; one brake B-1 which creates engagement of the fourth element; one brake B-2 which creates engagement of the second element; and one output element coupled with the third element R3.\nWith the scope of the Applicant's DE 101 15 983 A1 is described a multiple stage transmission having one drive shaft connected with a front-mounted range change set, one driven shaft connected with a rear-mounted range change set and having a maximum of seven shifting elements by the optional shifting of which at least seven forward gears can be shifting without group shift. The front-mounted range change set is formed by a front-mounted planetary gear set or a maximum of two non-shiftable front-mounted planetary gear sets coupled with the front-mounted planetary gear set, the rear-mounted range change set is designed as a two-spider-four shaft transmission with two shiftable rear-mounted range change sets and four free shafts. The first free shaft of said two-spider-four shaft transmission is connected with the first shifting element, the second free shaft with the second and third shifting elements, the third free shaft with the fourth and fifth shifting elements and the fourth free shaft is connected with the driven shaft. For a multiple stage transmission with a total of six shifting elements, the invention additionally proposes to connect the third free shaft or the first free shaft of the rear-mounted range change set with a sixth shifting element. For a multiple stage transmission having a total of seven shifting elements, the invention proposes connecting the third free shaft additionally with a shift shifting element D′ and the first free shaft additionally with a seventh shifting element.\nWithin the scope of the Applicant's DE 101 15 987 has been described a multiple stage transmission having at least seven gears. Said transmission, together with the drive shaft and the driven shaft, comprises one non-shiftable front-mounted range change set and one shiftable rear-mounted range change set in the form of a two-spider-four shaft transmission. The front-mounted range change set comprises one first planetary gear set which, together with the input rotational speed of the drive shaft, offers a second rotational speed which can optionally be shifted to a rear-mounted range change set. The rear-mounted range change set consists of two shiftable planetary gear sets which, with the six shifting elements, can shift at least seven gears, with two power paths being formed. At the same time, group shifts are always advantageously prevented during each shifting operation. One 9-gear multiple stage transmission has further become known from DE 29 36 969; it comprises eight shifting elements and four gear sets.\nAutomatically shiftable vehicle transmissions of a planetary design have already generally been widely described in the prior art and are subject to permanent developments and improvement. Said transmissions thus must have a sufficient number of forward gears, the same as one reverse gear and a ratio very well suited for motor vehicles with a high total spread and favorable ratio ranges. They must also make possible a high starting ratio in a forward direction and contain a direct gear, the same is adequate for use both in passenger and in commercial vehicles. In addition, said transmission must have low construction cost, especially require a low number of shifting elements and prevent double shifts in a sequential shifting mode so that in case of shiftings in defined gear groups only one shifting element is changed.\nThe problem on which this invention is based is to propose a multiple stage transmission of the type mentioned above in which the construction cost is optimized and, in addition, the efficiency degree is improved in the main drive gears with regard to towing and gearing losses. Besides, low torque must act on the shifting elements and planetary gear sets in the inventive multiple stage transmission, the same as the rotational speeds of the shafts, shifting elements and planetary gear sets are to be kept as low as possible. The number of gears, the same as the transmission spread must also be increased so that eight forward gears and at least one reverse gear can be advantageously implemented. Besides, the inventive transmission must be adequate for any design of the vehicle, especially for a front-transverse arrangement."}
{"text": "Systems may control access for entry to a facility or area based on data bearing records. Such systems may include the reading of data from data bearing records such as user cards or other articles or devices. Such data bearing records may be read through operation of a reader device, and the data read therefrom compared to stored data for purposes of determining if a bearer of the record is an individual that is authorized to access the facility or area. Such data bearing records may take different forms such as cards, tokens or other items from which data may be read. Such items may include magnetic snipe cards, radio frequency identification (RFID) cards, optically encoded cards, smart cards or other items from which data may be read.\nIn some embodiments additional data may be also read so as to identify a particular individual as authorized to access an area or facility. Such read data may be compared with stored data or correlated with other perceivable data to identify the individual as the appropriate authorized individual to be the possessor of the data bearing record with rights to access the facility. Such additional data may include the input of secret codes or personal identification numbers. Alternatively such data may include the sensing of biometric data, the receipt of voice data or other data that can be compared for purposes of identifying an authorized individual. Systems for controlling access may be used in connection with facilities where individuals who do not possess data bearing records such as a member of the public such as customers, visitors or the like, must also access the facility. Such individuals may also need to be provided access to or egress from the facility and the activities of such individuals controlled as may be appropriate.\nSuch systems may benefit from improvements."}
{"text": "In order to properly care for fish and other aquatic organisms contained within a reef aquarium, adequate circulation is required. The role of circulation is two fold: first, circulation acts to constantly mix the aquarium water itself, ensuring that proper chemistry is maintained throughout the entire aquarium. Adequate circulation maintains the equilibrium of oxygen and carbon dioxide by increasing the rate at which water flows from the bottom of the tank to the top, where it can take in these compounds from the air. The second role of circulation is related to the nature of the inhabitants of a reef aquarium. Because many reef inhabitants are sessile (they do not move), circulation is the only means by which nutrients such as food and oxygen are brought to these animals and the only means by which waste is expelled. In the ocean, corals and other sessile animals have the benefit of large waves crashing into the reef in a random but consistent fashion. Within the constraints of a glass box or aquarium, a pump is used as a substitute.\nPrior aquarium circulating devices and pumps feature two aspects that make them less ideal than the present invention. First, designs featuring epoxy sealed motors within the aquarium have the unfortunate side effect of being relatively large and distracting to the intrinsic beauty of an aquarium, add unwanted heat to the aquarium through direct contact with the motor stator, and require that electricity be brought into the aquarium itself via a power chord or a battery sealed into the motor assembly. Second, some prior designs utilize a mechanical bracket which hangs over the top of the aquarium in order to support the pump within the aquarium. In some prior pumps in which the motor and the centrifugal propeller are magnetically coupled through the glass, brackets are used to support and align the rotating component within the aquarium. The prior designs are unsatisfactory because they are bulky due to the motor being placed within the aquarium or due to the brackets supporting the motor outside the aquarium. Furthermore, the prior designs required that the pump be located at a location determined by the location of the bracket or be on the bottom of the aquarium due to the weight of the pump.\nThe present invention attempts to remedy these drawbacks and provides a fluid pump assembly adapted to be mounted to an aquarium without the use of mechanical aids, such as brackets. The disclosed pump can be located anywhere on the surfaces of the aquarium, thus maximizing the aesthetic effects of the aquarium and facilitating water circulation by allowing the pump to located at a location achieved optimized fluid flow based upon the interior characteristics of the aquarium."}
{"text": "A variety of systems have been proposed that call for the ability for wireless sensor networks (WSNs) to determine the location of a node within a wireless communication system. For example, in asset control, it is desirable to know the locations of objects (e.g., laptop computers, cell phones, shipping containers with a built-in sensor node, shared hospital equipment with either an attached or built-in sensor node, . . . , etc.) within the confines of, say, an office building, factory floor, or other general or specialized space. Unfortunately, many location applications for wireless sensor networks are characterized by their severe resource constraints (e.g. energy, bandwidth, processing power, and memory), which have a considerable impact on the location performance. For example, many WSN applications require very long lifetime to avoid frequent re-charging or battery replacement. A difficulty in designing an efficient location technique is managing the tradeoff between resource usage and the location performance. Minimizing resource consumption leads to degradation in location accuracy, while increasing resource consumption leads to a more accurate location estimate. Therefore a need exists for a method and apparatus for determining the range and location of a node within a wireless communication system that balances resource constraints with location accuracy.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions and/or relative positioning of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present invention. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are often not depicted in order to facilitate a less obstructed view of these various embodiments of the present invention. It will further be appreciated that certain actions and/or steps may be described or depicted in a particular order of occurrence while those skilled in the art will understand that such specificity with respect to sequence is not actually required. Those skilled in the art will further recognize that references to specific implementation embodiments such as “circuitry” may equally be accomplished via replacement with software instruction executions either on general purpose computing apparatus (e.g., CPU) or specialized processing apparatus (e.g., DSP). It will also be understood that the terms and expressions used herein have the ordinary technical meaning as is accorded to such terms and expressions by persons skilled in the technical field as set forth above except where different specific meanings have otherwise been set forth herein."}
{"text": "1. Field of the Invention\nThe present invention relates to a misfire discriminating method for an engine. More particularly, it relates to a method of discriminating the misfire of an engine for each cylinder wherein the misfire condition of each cylinder is discriminated from the momentum change rate between cylinders. Also, it relates to such a method wherein the misfire judgement level of each cylinder is set with reference to a revolution fluctuation arising when fuel has been cut off, and wherein the misfire condition of each individual cylinder is discriminated on the basis of the misfire judgement level.\n2. Description of the Prior Art\nIn general, it is ideal for producing a stable output that combustion in a multicylinder engine undergoes the same process every cycle. In the multicylinder engine, however, the combustion is liable to deviate due to the synergy of such drawbacks as listed below:\n(1) Nonuniformity in the distribution rate of intake air attributed to the complication of the shape of an intake pipe, the interferences of the intake air among cylinders, etc. PA1 (2) Some differences among the combustion temperatures of the individual cylinders attributed to cooling routes. PA1 (3) Manufactural dispersions in the volumes of the combustion chambers of the individual cylinders, the shapes of pistons, etc. PA1 (4) Slight discrepancies of the air fuel ratios of the individual cylinders caused by unequal fuel injection quantities ascribable to the manufactural errors of injectors, etc.\nHeretofore, the combustion fluctuations among the cylinders have been suppressed to the minimum by the air fuel ratio controls and ignition timing controls of the individual cylinders. In a recent high-performance engine tending toward a higher output and a lower fuel cost, however, when any of the injectors, ignition plugs etc. has degraded or broken down, intermittent misfire is liable to arise and to incur lowering in the output.\nEven when the intermittent misfire has occurred in one cylinder in the multicylinder engine, a driver often drives an automobile without noticing the misfire. Besides, it is difficult to judge during the drive whether the cause of the misfire is merely temporary or is the degradation or the like of any of the injectors, the ignition plugs etc.\nTherefore, according to the official gazette of Japanese Patent Application Laid-open No. 258955/1986 by way of example, a comparison is made between the difference of the minimum value and maximum value of the r. p. m. of an engine corresponding to a cylinder at the previous combustion stroke and the difference of the minimum value and maximum value of the engine r. p. m. corresponding to the cylinder at the present combustion stroke. The combustion condition of the pertinent cylinder is discriminated, depending upon whether or not the discrepancy between the compared values falls within a preset reference value. In a case where abnormal combustion has occurred in excess of a predetermined number of times, a misfire is judged, and warning is given.\nWith the prior-art technique, the combustion fluctuation of each cylinder is obtained from the the differences between the minimum engine r. p. m. and maximum engine r. p. m. values of the combustion stroke cylinder. During combustion, however, the engine speed rises abruptly, and the engine is loaded relatively heavily, so that the fluctuation of an acceleration increases. Accordingly, it is difficult to specify the maximum value of the engine r. p. m., and a precision error developing at the judgement of the misfire enlarges.\nMeanwhile, the combustion characteristics of an engine deviate, not only among cylinders, but also among individual engines on account of the manufactural errors of components, etc.\nWhen the reference value for comparing the revolution fluctuations is set as an absolute value as in the prior-art technique, it changes relatively every engine due to the deviation of the combustion characteristics of the individual engines, and it becomes difficult to precisely grasp abnormal combustion in some of the engines.\nIn an engine having a small number of cylinders, the combustion intervals between the cylinders are comparatively long, and hence, the difference of the revolution fluctuations is large. Therefore, even when the reference value is set as the absolute value, the deviation of the combustion characteristics of the individual engines do not exert great influence on the misfire judgement. In contrast, in an engine having a large number of cylinders, the combustion intervals are short, and the difference of the revolution fluctuations decreases to that extent. Therefore, when the judgement level (the reference value) is set as the absolute value beforehand, the deviation of the combustion characteristics of the individual engines affects the precision of the misfire judgement greatly.\nEspecially in a high speed region, the fluctuation difference becomes small. Therefore, when the judgement level changes relatively every engine, the precise misfire judgment becomes very difficult.\nBy way of example, with a technique in the official gazette of Japanese Patent Application Laid-open No. 82534,1984, a differentiated r. p. m. value .DELTA.N#i (.DELTA.N#i=NH#i-NL#i) which is the difference between the instantaneous engine r. p. m. NL#i of each of cylinders #i (i=1-4 in a four-cylinder engine) before a combustion stroke and the instantaneous engine r. p. m. NH#i thereof after the combustion stroke is obtained every cylinder, and the differentiated r. p. m. value .DELTA.N#i of each cylinder #i is subsequently compared with the average value .DELTA.NA of the differentiated r. p. m. values .DELTA.N#i of all the cylinders #i, thereby to grasp the combustion condition of each cylinder #i.\nWith the prior-art technique, however, a reference value is the average differentiated r. p. m. value .DELTA.NA of all the cylinders, which in itself is always prone to fluctuate depending upon the combustion conditions. When the average differentiated r. p. m. value .DELTA.NA of all the cylinders has fluctuated, that combustion condition of each cylinder which is estimated with reference to the average differentiated r. p. m. value .DELTA.NA of all the cylinders comes to involve the combustion condition factors of the other cylinders in the case of setting this value .DELTA.NA, and it becomes difficult to precisely grasp the misfire of each cylinder.\nMeanwhile, with a technique wherein as disclosed in the official gazette of Japanese Patent Application Laid-open No. 268956/1988, misfire is decided from the difference of momentums before and after combustion (a \"revolving circumferential difference\" in the official gazette), the misfire is simply judged on condition that the momental difference is minus. In, for example, the running of an automobile at a constant speed, however, the momental difference sometimes becomes minus under the influence of the friction of an engine in spite of normal firing (normal combustion), and the misfire might be erroneously decided."}
{"text": "Packaged semiconductor dies, including memory chips, microprocessor chips, and imager chips, can include semiconductor dies mounted on a package substrate. The semiconductor dies are encased in a plastic protective covering, and each die includes functional features, such as memory cells, processor circuits, and imager devices. Bond pads on the dies are electrically connected between the functional features and terminals on the package substrate that allow the dies to be connected to external circuitry.\nTo increase the density of dies within a package, the dies can be stacked upon one another within the casing. One challenge with vertically stacked dies, however, is that the dies can have different sizes or footprints. For example, in a memory package, a memory controller die can have a smaller footprint than the memory dies within the package. The memory controller die can be more difficult to wirebond because it is offset from the memory dies. Also, the memory dies can sometimes tilt when stacked upon the smaller memory controller die."}
{"text": "A conventional measurement device (e.g., a vehicle oscilloscope or a multimeter), especially one that does not have built-in wireless capability, requires a wired connection to a personal computer (PC) for controlling the measurement device and displaying various waveforms and measured values. A typical use of the conventional measurement device is to locate problems in a vehicle, following a troubleshooting tree or procedure. In most cases, the conventional measurement device has to be connected to circuits or components under the hood of the vehicle, while the PC and the operator stay in the driver's seat.\nUnfortunately, having such a wired connection of conventional measurement device to an external computer is unwieldy and not flexible (e.g., in the above-noted use scenario). In addition, the conventional measurement device has to be separately powered, which also adds to the unwieldiness of the conventional measurement device during measurements. These and other problems, which have been identified by the inventors, exist in conventional systems. Thus there is a need for a wireless adapter for a measurement device."}
{"text": "A. Field of Invention\nThis invention pertains to the art of methods and apparatuses related to vehicle construction and more specifically to methods and apparatuses related to the connection of vehicle roof panels to side panels.\nB. Description of the Related Art\nIt is well known in the automotive industry to attach roof panels to side panels. Typically, the roof panel is formed of steel and the side panel (sometimes referred to as a side panel outer, or SPO) is also formed of steel. In this case, it is relatively easy to connect the roof panel to the side panel in a welding operation.\nIt is also known, however, to form roof panels out of various thermoplastics, such as fiber reinforced (composite) plastics. This usually means the roof panel is formed of a carbon fiber reinforced thermoplastic, or the like. In some specific cases, the thermoplastic roof panel may be transparent. A problem with thermoplastic roof panels is that they are more difficult to connect to the steel (or other metal) side panels because the different materials make the panels incompatible for welding operations.\nVarious methods for connecting thermoplastic roof panels to steel side panels are known. They include the use of adhesives, bolts, and rivets. While these methods and apparatuses generally work well for their intended purposes, they have disadvantages. Disadvantages to the use of adhesives include relatively high cost and relatively limited connection strength. Disadvantages to the use of bolts include increased weight and increased manpower for assembly. Disadvantages to the use of rivets include limitations in types of rivets that can be used, relatively high cost, relatively limited connection strength, and increased manpower for assembly. A common disadvantage to all types of fasteners (whether bolts, rivets, or other fasteners) is that their visibility is often considered undesirable to customers.\nTherefore, what is needed is a method for connecting thermoplastic roof panels to side panels in a manner to reduce or eliminate the disadvantages known in the art."}
{"text": "1. Field of the Invention\nThe present invention relates to an air conditioner.\n2. Description of Related Art\nJP-A-2000-255257 describes an air conditioner for a vehicle which uses photocatalyst in a heat exchanger. A blower circulates air in a case of the air conditioner, and an evaporator is arranged in the case so as to cool the air. A deodorization filter is provided upstream of the evaporator in a flow of the air. The deodorization filter has adsorbent and photocatalyst. Further, a surface of the evaporator has photocatalyst. Moreover, an ultraviolet ray lamp is disposed between the evaporator and the deodorization filter.\nThe ultraviolet ray lamp is turned on always or periodically with interval during operation of the blower. The photocatalyst of the deodorization filter is excited by the ultraviolet rays emitted from the lamp toward the deodorization filter, so that odorant adsorbed on the adsorbent is decomposed. The photocatalyst of the evaporator is excited by the ultraviolet rays emitted from the lamp toward the evaporator, so that bacteria on the surface of the evaporator is sterilized.\nHowever, because the ultraviolet ray lamp is arranged at a middle of the case, ventilation air resistance of the air circulating through the case is increased, so that air sending amount of the blower is decreased."}
{"text": "1. The Field of the Invention\nThe present invention relates to a means for mounting a magnetic read head in a card reader or the like.\n2. The Prior Art\nThere have been many successful schemes for mounting magnetic read heads in card readers and the like. These mounting assemblies have usually included a gimbaling means which allow the magnetic read head to have a substantial amount of float to accommodate irregularities in the cards. However, the previous mounting means have all been somewhat cumbersome in that they have required a large number of parts and a large amount of space in order to arrive at an arrangement which will accommodate movement of the magnetic read head about several axes. An example of a prior art mounting means can be found in U.S. Pat. No. 4,167,664."}
{"text": "The expression of genes during the development of a pluripotent or progenitor cell into a differentiated, mature cell can provide a context for the study of tumorigenic cells whose origin is derived from such progenitor cells. In certain hematopoietic or epithelial tumors, malignant gene expression correlates substantially with the expression observed during normal development of the tissue from which the tumor originates, Gordon et al., J. Cell Biol. 108: 1187 (1989); Godal et al., Adv. Cancer Res. 36: 211 (1982). In fact, many biological activities of progenitor cells, including cellular migration and tissue remodeling, resemble pathological activities of cancer cells, such as metastases and tumor invasion.\nNeuroblastoma, a tumor of the adrenal gland which afflicts persons during early childhood, is another system in which tumor biology correlates with that of normal differentiation and morphogenesis of its progenitor cells (neuroblast). Neuroblastoma is an embryonal tumor that exhibits both undifferentiated and differentiated histopathology. The development of neuroblastoma tumors mimics stages identifiable during histogenesis of its tissue of origin, the adrenal medulla. Cooper et al., Cell Growth and Diff. 1: 149 (1989).\nDuring the development of human adrenal medulla neuroblast into mature chromaffin cells, four individual genes are expressed in a sequential pattern. Once a neuroblast is induced to differentiate along a neuroendocrine pathway, the progressive stages of chromaffin maturation are marked by a temporal expression of genes denoted TH, CGA, pG2 and B2M (Cooper, supra. at page 153). Cooper identified that the pattern of gene expression of these four markers in neuroblastoma cells mimics that of normal adrenal neuroblast arrested during three different stages of development.\nOne of these marker genes, pG2, was identified first in pheochromocytoma, a tumor of the adult adrenal medulla (Helman et al., PNAS USA 84: 2336 (1987)). Helman reported that pG2 also is highly expressed normal human adrenal cells.\nHelman isolated a full-length cDNA from a human adrenal cDNA library, and identified a corresponding pG2 protein containing 286 amino acids, having a predicted molecular weight of 30,600 daltons (Helman et al., Nucleic Acids Res. 18(3): 685 (1990)).\nA gene having developmentally-regulated expression, paralleling that of pG2, would be useful for detecting pheochromocytoma or neuroblastoma by genetic methods, especially since pG2 expression is restricted to the adrenal gland in non-malignant tissues."}
{"text": "1. Field of the Invention\nThis invention relates generally to a balancer for dehydration tubs for use in washing machines or the like, and more particularly to such a balancer of the liquid-in-container type including a concentric multi-compartment container.\n2. Description of the Prior Art\nIn dehydration tubs used in washing machines having a centrifugal dehydrating function or the like, a liquid-in-container type balancer is mounted on an upper open end of the dehydration tub for the purpose of correcting an unbalanced condition thereof due to one-sided laundry. For improvement in a correcting force, the interior of the balancer is divided by one or more vertically extending concentric partition walls into a plurality of compartments. A predetermined amount of liquid (usually, salt water) is contained in each of the compartments.\nFIGS. 23 and 24 illustrate the construction of a conventional balancer of the type described above. Referring to FIG. 24, the balancer comprises an annular balancer container 1 having in its interior a partition wall 2 standing from the bottom thereof. The interior of the container 1 is divided by the partition wall 2 into inner and outer compartments 3 and 4. A lid 5 is attached to the top of the container 1 so as to close upper openings of the compartments 3 and 4. The lid 5 has two inlets 6 and 7 formed therein to correspond to the compartments 3 and 4 respectively. A liquid 8 is poured through the inlets 6 and 7 into the respective compartments 3 and 4. Thereafter, closures 9 and 10 are closely fitted into the inlets 6 and 7 respectively.\nIn the balancer described above, the liquid 8 contained in the container 1 leaks out or flows between the compartments 3 and 4 if the compartments are not watertightly sealed by the lid 5 or if the inlets 6 and 7 are not watertightly closed by the respective closures 9 and 10. Consequently, an expected correcting force cannot be obtained when the liquid 8 leaks out of the container 1 or flows between the compartments 3 and 4. In view of this problem, watertight tests need to be carried out for the balancer.\nThe watertight tests are carried out in the following procedure. The interior of the compartment 3 is pressurized or depressurized through the inlet 6 after the lid 5 has been attached to the container 1. Consequently, the watertightness is tested at a boundary 11 between an inner wall of the container 1 and the lid 5 and at a boundary 12 between the partition wall 2 and the lid 5. Subsequently, the interior of the compartment 4 is pressurized or depressurized through, the inlet 7 in order that the watertightness is tested at a boundary 13 between an outer wall of the container 1 and the lid 5. Finally, the inlets 6 and 7 are closed by the respective closures 9 and 10 after the liquid 8 has been poured into the compartments 3 and 4 through the inlets respectively. The completed balancer is then put into a chamber, and the interior of the chamber is pressurized or depressurized in order that the watertightness at boundaries between the circumferential edges of the inlets 6 and 7 and the closures 9 and 10 is tested.\nIn the conventional balancer, the watertight test needs to be carried out for each of the compartments 3 and 4, and furthermore, the other watertight test needs to be carried out at a stage of the end product of balancer. Accordingly, steps of the watertight test is increased with an increase in the number of compartments of the container. Furthermore, since the portions of the balancer to be tested are diverse, an equipment for the watertight test is rendered complicate, and accordingly, the cost of equipment is increased.\nOn the other hand, the number of compartments of the container or the height of the balancer is increased in the prior art so that the performance of the balancer is improved. In each case, the size of the balancer is increased, which results in an increase in the size of the washing machine or a decrease in a washing capacity of the washing machine.\nFurthermore, the correcting performance of the balancer mainly depends upon the configuration of the container 1. The correcting performance of the balancer cannot be altered after its configuration has been determined."}
{"text": "Machine learning and classification involves the extraction of structured or unstructured targeted raw data from various data sources and transforming the targeted raw data into a format suitable for performing subsequent learning and/or classification processing on. During machine learning, descriptive patterns within the target data are collected and documented using a consistently applied aggregation method, which typically collects descriptive statistical facts about data of known categories or classification origins. When data of unknown categories or classification origins are used as data input, the consistently applied aggregation method can be used to identify descriptive patterns within the unknown data for comparison and classification against all data of known classification types and origins within the current system.\nA parallel data processing implementation for the purposes of machine learning and classification typically comprises the distribution of targeted raw data and a consistently applied aggregation method across multiple processors and disks to perform simultaneous computations, while taking advantage of increased memory, reduced i/o, and the combined processing power of all participating computing devices.\nSeveral systems and programming languages, such as the R Project for Statistical Computing, include programming features for the learning and classification of data. However, current solutions do not provide the functionality to automatically distribute learning and classification processes across multiple processors and disks in a distributed parallel computing environment using a map reduction aggregation method. Furthermore, current learning and classification systems implement a rigid framework which typically requires the use of a single predefined aggregation method and classification metric function. These limitations often result in extensive and sometimes very overhead intensive input data pre-processing in order to transform the targeted data for use within the rigid framework. Furthermore, the rigid framework typically does not support multiple application specific map reduction aggregation methods and classification metric functions created by the application programmer."}
{"text": "1. Field of the Invention\nThe present invention relates to an energy recovery circuit and energy recovering method using the same, and more particularly, to an energy recovery circuit and energy recovering method using the same that is capable of reducing the number of components.\n2. Description of the Related Art\nRecently, there have been developed various flat panel display devices reduced in weight and bulk that is capable of eliminating disadvantages of a cathode ray tube (CRT). Such flat panel display devices include a liquid crystal display (LCD), a field emission display (FED), a plasma display panel (PDP) and an electro-luminescence (EL) display, etc.\nThe PDP among them is a display device using gas discharge and has an advantage that it can be easily produced in a large sized panel. As shown in FIG. 1, a three electrode AC surface discharge PDP is typical as the PDP, wherein it has three electrodes and is driven by AC voltage.\nReferring to FIG. 1, a discharge cell of a three-electrode, AC surface-discharge PDP includes a scan electrode 12Y and a sustain electrode 12Z provided on an upper substrate 10, and an address electrode 20X provided on a lower substrate 18.\nOn the upper substrate 10 provided, in parallel, with the scan electrode 12Y and the sustain electrode 12Z, an upper dielectric layer 14 and a protective film 16 are disposed. Wall charges generated upon plasma discharge are accumulated onto the upper dielectric layer 14. The protective film 16 prevents a damage of the upper dielectric layer 14 caused by a sputtering during the plasma discharge and improves the emission efficiency of secondary electrons. This protective film 16 is usually made of magnesium oxide (MgO).\nA lower dielectric layer 22 and barrier ribs 24 are formed on the lower substrate 18 provided with the address electrode 20X. The surfaces of the lower dielectric layer 22 and the barrier ribs 24 are coated with a phosphorous material 26. The address electrode 20X is formed in a direction crossing the scan electrode 12Y and the sustain electrode 12Z. The barrier rib 24 is formed in parallel to the address electrode 20X to thereby prevent an ultraviolet ray and a visible light generated by a discharge from being leaked to the adjacent discharge cells.\nThe phosphorous material 26 is excited by an ultraviolet ray generated during the plasma discharge to generate any one of red, green and blue visible light rays. An inactive mixture gas for a gas discharge is injected into a discharge space defined between the upper and lower substrates 10 and 18 and the barrier rib 24.\nThe three electrode AC surface discharge PDP is divided into a plurality of subfields to be driven, wherein the light emission of numbers proportional to the weight of a video data is in progress in each subfield period, thereby performing the gray level display. The subfield is re-divided into an initialization period, an address period, a sustain period and an erasure period to be driven.\nHerein, the initialization period is a period when uniform wall charges are formed in a discharge cell, the address period is a period when a selective address discharge is generated in accordance with the logical value of the video data, the sustain period is a period when a discharge is kept in the discharge cell where the address discharge is generated, and the erasure period is a period when the sustain discharge generated during the sustain period is eliminated.\nIn AC surface discharge PDP driven like this way, a high voltage of not less than several hundreds of volts is required in the address discharge and the sustain discharge thereof. Accordingly, an energy recovery circuit is used for minimizing a drive power required in the address discharge and the sustain discharge. The energy recovery circuit recovers the voltage between the scan electrode 12Y and the sustain electrode 12Z, and utilizes the recovered voltage as a drive voltage for the next discharge.\nFIG. 2 is a diagram illustrating an energy recovery circuit installed for recovering a voltage of the sustain discharge.\nReferring to FIG. 2, energy recovery circuits 30, 32 of the related art PDP are symmetrically installed with a panel capacitor Cp, therebetween. Herein, the panel capacitor Cp equivalently represents the capacitance which is formed between the scan electrode Y and the sustain electrode Z. In the energy recovery circuits, a first energy recovery circuit 30 supplies a sustain voltage to the scan electrode Y and a second energy recovery circuit 32 supplies the sustain voltage to the sustain electrode Z while it alternately operates with the first energy recovery circuit 30.\nThe composition of the energy recovery circuits 30, 32 of the related art PDP is described in reference with the first energy recovery circuit 30. The first energy recovery circuit 30 includes an inductor L connected between a panel capacitor Cp and a source capacitor Cs; first and third switches S1, S3 connected in parallel between the source capacitor Cs and the inductor L; and second and fourth switches S2, S4 connected in parallel between the panel capacitor Cp and the inductor L.\nThe second switch S2 is connected to a sustain voltage source Vs, and the fourth switch S4 is connected to a ground voltage source GND. The source capacitor Cs recovers the voltage charged into the panel capacitor upon the sustain discharge to be charged and re-supplies the charged voltage to the panel capacitor Cp. The voltage of Vs/2 corresponding to the half value of the sustain voltage source Vs is charged in the source capacitor Cs. The inductor L forms a resonance circuit together with the panel capacitor Cp. For this, the first to fourth switches S1 to S4 control the flow of electric current.\nOn the other hand, fifth and sixth diodes D5, D6 each installed between the first and third switches S1, S3 and the inductor L prevent the current from flowing in a reverse direction.\nFIG. 3 is a timing diagram and waveform diagram representing an output waveform of a panel capacitor and an on/off timing of switches of the first energy recovery circuit.\nBefore a T1 period, assuming that a voltage of 0 volt is charged in the panel capacitor Cp and a voltage of Vs/2 is charged in the source capacitor Cs, the operation process is described in detail.\nIn the T1 period, a first switch S1 is turned on to form a current path from the source capacitor Cs to the panel capacitor Cp through the first switch S1 and the inductor L. Accordingly, the voltage of Vs/2 charged in the source capacitor Cs is supplied to the panel capacitor Cp. At this moment, the inductor L and the panel capacitor Cp forms a series resonance circuit, thus the sustain voltage Vs which is double of the voltage of the source capacitor Cs is charged in the panel capacitor Cp.\nIn a T2 period, the second switch S2 is turned on. When the second switch S2 is turned on, the voltage from the sustain voltage source Vs is supplied to the scan electrode Y. The voltage of the sustain voltage source Vs supplied to the scan electrode Y prevents the voltage of the panel capacitor Cp from dropping below the sustain voltage source Vs to cause the sustain discharge to be generated in a normal manner. On the other hand, the voltage of the panel capacitor Cp rises to the sustain voltage Vs in the t1 period, thus the drive power supplied from the outside to generated the sustain discharge is minimized.\nIn a T3 period, the first switch S1 is turned off. At this moment, the scan electrode Y maintains the voltage of the sustain voltage source Vs for the T3 period. In a T4 period, the second switch S2 is turned off and the third switch is turned on. When the third switch S3 is turned on, there is formed a current path from the panel capacitor Cp to the source capacitor Cs through the inductor L and the third switch S3 to recover the voltage charged in the panel capacitor Cp to the source capacitor Cs. At this moment, the source capacitor Cs is charged with the voltage of Vs/2.\nIn a T5 period, the third switch S3 is turned off and the fourth switch S4 is turned on. When the fourth switch S4 is turned on, a current path is formed between the panel capacitor Cp and the ground voltage source GND, thus the voltage of the panel capacitor Cp drops to 0V. In a T6 period, it maintains at the T5 state for a designated period. In fact, an AC drive pulse supplied to the scan electrode Y and the sustain electrode Z is obtained while the T1 to T6 periods are repeated periodically.\nOn the other hand, as shown in FIG. 4, the second energy recovery circuit 32 supplies the drive voltage to the panel capacitor Cp while alternately operating with the first energy recovery circuit 30. Accordingly, the panel capacitor Cp receives the sustain pulse voltage Vs that has a different polarity as shown in FIG. 4. In this way, the sustain pulse voltage Vs having the different polarities is supplied to the panel capacitor Cp, thus the sustain discharge is generated at the discharge cells.\nHowever, since the first energy recovery circuit 30, installed at a side of the scan electrode Y, and the second energy recovery circuit 32, installed at a side of the sustain electrode Z, are respectively operated, lots of circuit components such as switching device are required. Accordingly, there is a problem that a manufacturing cost thereof becomes increased. In addition, if lots of circuit components are installed to the energy recovery circuits 30, 32, then a large amount of power consumption becomes wasted."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for decomposing perfluorinated compounds and to a system for processing perfluorinated compounds.\nA claim of priority is made to Korean Patent Application No. 2002-23508, filed in the Korean Intellectual Property Office on Apr. 29, 2002, the entirety of which is incorporated herein by reference.\n2. Description of the Related Art\nPerfluorinated compounds (PFCs) are toxic waste gases generated and exhausted into the air during the manufacture of semiconductor devices. PFCs are commonly used in dry etching, chemical vapor deposition, and chamber cleaning processes. Typical examples of PFCs include CF4, C2F6, C3F8, CHF3, NF3, SF6, and the like. PFCs are very stable compounds inherently having a strong binding force, so they are resistant to decomposition and have a long lifespan. Furthermore, the global warming potential of PFCs is much higher than that of carbon oxide, which has prompted the Word Semiconductor Council to voluntarily take steps to reduce the discharge of the PFCs.\nA conventional technique for processing PFCs exhausted during the manufacture of semiconductor devices is disclosed in U.S. Statutory Invention Registration (SIR) No. H1701. According to this method, fluorine in the exhaust gas is reacted with aluminum to produce AlF3 as a waste form. However, the addition of aluminum increases costs. In addition, the requirement for additional equipment for processing the reaction product increases the overall processing system size and maintenance and management functions.\nThere exist a variety of PFC processing methods, including a method involving the addition of alkaline earth metals and a powerful reducing agent, a method involving high-temperature combustion, and a method involving the removal of particles from the exhaust gas, the addition of hydrogen or water and oxygen, and thermal decomposition above 600° C. Recently, a method for removing PFCs through exposure to a process gas in a plasma state at 10,000K and rapid cooling has been proposed. However, in such a high-temperature thermal plasma technique, it is necessary to head the process gas to at least 600° C. and to keep this high temperature level for period of time, and thus energy consumption is high. Moreover, the throughput and performance are not high when compared with labor and cost inputs.\nAccording to most techniques using low-temperature plasma, similar to the above-described U.S. SIR, PFCs are converted into solid waste particles by the addition of alkaline earth metals or other additives. However, as described above, these techniques are considered to be uneconomical due to the need for collateral equipment, and the accompanying maintenance and management costs."}
{"text": "This type of photocatalysts are finding expanded applications because energy used in catalytic reactions is light energy, such as sunlight, that is inexhaustible. For example, titanium dioxide (TiO2), a kind of photocatalyst, particularly titanium oxide in an anatase crystal form, produces excited electrons and positive holes upon exposure to light energy (ultraviolet light), and the excited electrons and positive holes produce active oxygen species, such as O2−, O−, and •OH (wherein the dot • represents an unpaired electron and means that the species attended with the dot • is a free radical), in the presence of oxygen and water on the surface of the catalyst. Applications utilizing free radical properties of the active oxygen species have been proposed such as air purification applications where nitrogen oxides (NOx) in the air are oxidized with the active oxygen species and consequently converted to a harmless reaction product (nitric acid), and degradation of bacteria through the oxidation of organic matter, that is, the so called antimicrobial applications.\nIn the course of the oxidation of the nitrogen oxides with the active oxygen species, nitrogen dioxide (NO2) is produced as an intermediate that is further oxidized and finally converted to nitric acid. As a result of the production of nitric acid, the nitrogen oxides in the air are reduced, and the air is purified. For this reason, the presence of the active oxygen species and the nitrogen oxide or nitrogen dioxide is indispensable for enhancing the percentage reduction of nitrogen oxides. Since, however, nitrogen dioxide is a relatively chemically stable compound (gas), the produced nitrogen dioxide leaves the reaction system. This lowers the efficiency of the oxidation with the active oxygen species, resulting in lowered percentage reduction of the nitrogen oxides. Use of porous adsorbents, such as activated carbon, is considered effective for preventing nitrogen dioxide from leaving the reaction system. As is apparent from the following description, this method is not always effective.\nSpecifically, when nitrogen dioxide, which has left, is adsorbed on the above adsorbent, the nitrogen dioxide often remains within pores of the adsorbent without being released. For this reason, in some cases, the adsorbed nitrogen dioxide is placed outside the system of oxidation with the active oxygen species and does not undergo the oxidation reaction and hence cannot be converted to nitric acid as a final product. Thus, the nitrogen oxides are not finally converted to nitric acid. This inhibits the reduction of nitrogen oxides. In this case, it should be noted that nitrogen dioxide adsorbed onto the adsorbent in a region where nitrogen dioxide can be present together with the active oxygen species and is in the reaction system, that is, in a region close to the photocatalyst, is oxidized to produce nitric acid. Since, however, the region close to the photocatalyst occupies only a small proportion of the whole material adsorption region (including pores) in the absorbent, it can be said that the proportion of nitrogen dioxide, which could not be oxidized to nitric acid, is high. That is, the adsorbent merely adsorbs and holds nitrogen dioxide, and the percentage reduction of nitrogen oxides by conversion to nitric acid does not appear to be satisfactory.\nThe present invention has been made with a view to solving the above problems, and an object of the present invention is to further improve the efficiency of a catalytic reaction, in which a photocatalyst participates, or to improve the percentage reduction of the reactant applied to the catalytic reaction through the conversion of the reactants to the final product. Another object of the present invention is to supplement the function of a photocatalyst."}
{"text": "A wide variety of devices have been used in the past for controlling the speed of an electric motor ranging from a simple variable resistor connected in series between the motor and a power source to rather sophisticated control devices. The commonly used variable resistor is not suitable for use with a shaded pole motor, often used with laboratory equipment, because it reduces the starting torque such that the motor may not start at low speeds.\nOne known method for controlling motor speed is to switch a bidirectional semiconductor device, such as a triac, on and off to control the delivery of the power source to the motor. See, for example, U.S. Pat. No. 3,780,366 to Henderson; U.S. Pat. No. 3,523,234 to Turtle; U.S. Pat. No. 3,541,412 to Worth; U.S. Pat. No. 3,596,158 to Waltour; U.S. Pat. No. 3,678,360 to Minarik et al.; U.S. Pat. No. 3,742,337 to Digneffe; U.S. Pat. No. 3,803,468 to Soeda; and U.S. Pat. Nos. 3,875,485 and 3,848,167 to Hornung. Such devices could be used with shaded pole motors; however, the devices shown in these patents are either very complicated and expensive or utilized some feedback mechanism coupled to the motor speed for adjusting the speed of a motor and are just not practical for replacing simple variable resistor controls.\nAccordingly, it is an object of the present invention to make a motor speed control circuit which utilizes a bidirectional switch to control the power supplied to the motor in a simple, inexpensive manner and which is suitable for use with shaded pole motors.\nIt is also known in the art that the speed of a motor is directly related to the number of energy pulses supplied to it over a period of time. See, U.S. Patent No. 4,104,570 to Hamby et al; U.S. Pat. No. 3,243,677 to Cannalte et al.; and U.S. Pat. No. 3,559,017 to Dinger. Cannalte et al. switch a D.C. voltage source on and off in a controlled manner to vary the speed of a motor. Dinger utilizes a complicated counting circuit to supply the proper number of pulses to the motor which matches the number of pulses which should be supplied to obtain a certain motor speed. Hamby et al. supply a number of pulses to six thyristors to adjust the motor speed or phase of firing signals from desired values. These devices are not suitable for replacing the variable resistor used to control shaded pole motors.\nTherefore, it is a further object of this invention to provide for a device which controls the speed of a motor directly by supplying a specific number of complete A.C. signal cycles to the motor in a simple and inexpensive manner."}
{"text": "It is well-known to provide in an office or work environment for one or a plurality of workstations to be used by one or a plurality of workers. In such known workstations, it is well-known to provide one or more worksurfaces as well as one or more storage units (e.g. drawers, cabinets, shelves, etc.). A typical workstation may include other devices or apparatus that facilitate the work of the worker (e.g. the use and communication of information), such as computing devices, task or other lighting, telephony, as well as connections for utilities (e.g. power, voice, data, etc.).\nIt is also known to provide a “threshold” (or “utility threshold”) having a generally vertical base frame and a generally vertical overhead top portion for use within a workstation, for example, as shown in U.S. Pat. No. 6,374,547. According to such known arrangements, the threshold is movable within the workstation along a path of travel so that it may be used to “divide” the work space and/or to provide privacy for the work of one or more workers within the workstation; the movable threshold may also be configured to deliver utilities within the workstation.\nIt is further known to provide for a worksurface that is movable within a workspace, so that one or more workers or other persons (who may need to work or collaborate in a wide variety of conditions) may place the worksurface in a position that is more beneficial for a particular purpose. According to certain known arrangements, the movable worksurface may be provided with one or more information display devices (e.g. panel displays or monitors of the type associated with a computing device such as a personal computer), for example, as shown in U.S. patent application Ser. No. 09/887,519.\nIt is also known to provide for a mobile worksurface, such as a mobile table, in a work space. In a common application, the mobile worksurface (alone or in conjunction with fixed worksurfaces), enhances the utility of a work space by providing greater flexibility and a greater range of use and other conditions. However, mobile worksurfaces—whether in use or not in use—typically require at least a portion of the work space, and at times may prove to be obstacles to occupants of the work space. The range of motion typically provided by such mobile worksurfaces is limited, for example to a particular plane or path of travel. In addition, mobile worksurfaces, typically provided with wheels or otherwise coupled to a fixed mounting structure, tend to lack the stability, size or “structure” typically provided with fixed worksurfaces. Moreover, mobile worksurfaces tend not to readily accommodate computing devices and/or display devices or other appliances and their associated cables (e.g. for power and/or data).\nAccordingly, it would be advantageous to provide for a workstation providing for an apparatus such as a threshold and/or worksurface that is conveniently movable from one position to another position. It would also be advantageous to provide for an apparatus such as a threshold and/or worksurface that provides a range of movement to cover a substantial portion of the work space within the workstation. It would further be advantageous to provide for a support system for an apparatus such as a threshold and/or worksurface that allows for movement within a horizontal plane or other paths of travel within the workstation. It would also be advantageous to provide for an apparatus such as a threshold and/or worksurface movable within a workstation that conveniently allows for association with a display device and thereby for convenient repositioning of the display device within a work space. It would also be advantageous to provide an apparatus such as a threshold and/or worksurface that conveniently provides for a wide range of motion and allows for a variety of orientations. It would further be advantageous to provide for an apparatus that conveniently provides for management and interconnection of cables providing utilities to the appliances or equipment on a worksurface or to a display device (or display devices) and/or associated with a threshold. It would further be advantageous to provide for an apparatus that can readily be integrated with the articles of furniture within a workstation.\nIt would be desirable to provide a workstation having any one or more of these or other advantageous features."}
{"text": "1. Field of the Invention\nA subject matter of the present invention is a cleaning and disinfecting composition of particular use in the cleaning and disinfecting of pipes and vessels of plants in the food industry, in particular the brewing industry, dairy industry, and the like.\n2. Description of the Related Art\nCompositions for cleaning and disinfecting plants in the farm produce industries are mainly used in systems known as xe2x80x9ccleaning in placexe2x80x9d (CIP) systems, according to which the plant is cleaned and disinfected without being disassembled, the disinfecting cleaning composition being circulated in the plant by pumping and/or applied by spraying (by means of permanently installed sprayers), so as to reach all parts of the plant to be treated.\nIn order to be applicable in a truly satisfactory manner in the food industry and more particularly in the brewing industry, such a disinfecting composition must combine several properties and must meet several requirements which it is difficult to satisfy simultaneously.\nA solution for use prepared by dilution of such a composition must in particular have a satisfactory microbiocidal action, both with regard to gram+ and gramxe2x88x92 bacteria and with regard to yeasts, such as Saccharomyces, Candida, and the like. A satisfactory microbiocidal effect is considered to have been obtained when the initial population of microorganisms is reduced by a factor of 105 in less than 20 minutes at 3xc2x0 C.\nSuch a solution for use must also be non-foaming or, in any case, exhibit very little foaming. This is because the formation of foam can result in cavitation of the pumps of the CIP system, causing non-circulation of the disinfecting solution and possibly destroying the pumps. In addition, the formation, even modest, of foam can result in difficulties during rinsing.\nSuch a solution for use must also have a good cleaning and optionally descaling power and be fully rinsable, that is to say without risk of specific adsorption on the surfaces of the materials commonly encountered in plants in the food industry.\nThe concentrated composition must be stable and homogeneous.\nThe solution for use prepared by dilution of this concentrated composition must also itself be homogeneous and sufficiently stable to allow it be recovered. This is because the solutions for use used in a CIP system are recovered after use and collected in storage tanks in order to be reused during a subsequent cleaning and disinfection operation on the plant.\nBefore a fresh use of a recovered solution for use, the concentration of this solution should generally be readjusted by the addition of a certain amount of the concentrated composition. It is consequently important for the concentration of the solution for use to be able to be easily measured continuously. When the solution for use comprises a constituent such as a strong acid, the concentration of active constituents in the solution is very easily measured by measuring the electrical conductivity of the solution, since this conductivity is approximately proportional to the concentration of strong acid, which concentration is proportional to that of the other active constituents of the solution for use.\nNumerous compositions for cleaning and/or disinfecting plants in the agricultural and food industries have already been provided.\nEP-A-0,147,102 discloses an antimicrobial composition comprising an inorganic acid and/or an organic acid, an xcex1-halogenated acetic acid and a linear or branched aliphatic fatty acid.\nAn antimicrobial composition which essentially comprises an aliphatic alcohol having from 6 to 12 carbon atoms and an aliphatic acid having from 2 to 6 carbon atoms, which acid is substituted at the xcex1-position by a chloro, iodo, cyano, nitro, amino, imino, thiohydroxyl or hydroxyl radical, is known from EP-A-0,208,403. This composition can also comprise an organic or inorganic acid.\nEP-A-0,524,075 discloses a disinfecting acid composition based on phosphoric acid and optionally on hydroxyacetic acid comprising a specific nonionic surface-active agent and an anionic hydrotropic agent of the arylsulfonate or alkylsulfonate type.\nA dilutable disinfecting cleaning composition comprising a fatty acid, a hydrotropic agent of the type of surface-active sulfonates, a strong acid and a stabilizing agent consisting of propionic acid, butyric acid, valeric acid or a mixture of these acids, is known from WO 94/10837.\nDisinfecting and cleaning compositions for the food industry are known from the document EP-A-0,245,928, these compositions comprising a monocarboxylic or dicarboxylic acid, a solubilizing agent, a diluent and an acid or a mixture of acids, in order to ensure a pH of less than 5. As specified in the document EP-B-0,245,928, the solubilizing agent must be an agent of the alkylated N,N-dimethylamine oxide type with 8 to 10 carbon atoms in the alkyl part in order for the composition to be stable and not very foaming.\nIt has now been found that compositions not comprising this solubilizing agent are stable and not very foaming and have an excellent effectiveness in combating bacteria and yeast. Furthermore, the compositions according to the invention can be easily discharged from the vessels and pipes by rinsing, so that any risk of detrimental change in the taste or flavor of the foodstuff by the presence of traces of compositions according to the invention can be avoided. The compositions according to the invention are therefore fungicidal and bactericidal compositions of low toxicity.\nThe disinfecting composition according to the invention, which is provided either in concentrated form or in dilute form, is a liquid composition comprising:\n(a) a monocarboxylic acid of formula\nRxe2x80x94COOH\nxe2x80x83with R a saturated or unsaturated, straight- or branched-chain alkyl radical comprising from 6 to 12 carbon atoms (preferably from 8 to 10 carbon atoms), or a mixture of such carboxylic acids;\n(b) a strong inorganic acid,\n(c) an organic acid of formula:\nRxe2x80x2CH2xe2x80x94CO2H\nxe2x80x83with Rxe2x80x2: H or OH,\nxe2x80x83or a mixture of such organic acids;\n(d) one or more anionic surfactants;\n(e) a solubilizing agent consisting of a compound of formula:\nRxe2x80x3O(CH2CH2O)xxe2x80x94H\nxe2x80x83in which\nRxe2x80x3 represents a C2 to C6 alkyl radical,\nx represents an integer from 1 to 3, or a mixture of such compounds;\n(f) water.\nThe composition according to the invention is a bactericidal and fungicidal composition.\nAccording to one embodiment of the concentrated composition, it comprises:\n(a) from 0.25 to 20% by weight of a monocarboxylic acid of formula:\nRxe2x80x94COOH\nxe2x80x83with R a saturated or unsaturated, straight- or branched-chain alkyl radical comprising from 6 to 12 carbon atoms, or of a mixture of such carboxylic acids;\n(b) from 20 to 40% by weight of a strong inorganic acid;\n(c) from 0.25 to 7% by weight of an organic acid of formula:\nRxe2x80x2CH2xe2x80x94CO2H\nxe2x80x83with Rxe2x80x2: H or OH,\nxe2x80x83or from 0.5 to 12% by weight of a mixture of such organic acids;\n(d) from 5 to 20% by weight of anionic surfactant(s), and\n(e) from 1 to 10% by weight of a solubilizing agent consisting of a compound of formula:\nRxe2x80x3O(CH2CH2O)xxe2x80x94H\nxe2x80x83in which\nRxe2x80x3 represents a C3 to C5 alkyl radical,\nx represents an integer from 1 to 3, or a mixture of such compounds.\nIn a particularly advantageous way, the composition comprises:\n(a) from 0.5 to 15% by weight of a monocarboxylic acid of formula:\nRxe2x80x94COOH\nxe2x80x83with R a saturated or unsaturated, straight- or branched-chain alkyl radical comprising from 6 to 12 carbon atoms, or of a mixture of such carboxylic acids;\n(b) from 25 to 35% by weight of a strong inorganic acid;\n(c) from 2 to 5% by weight of an organic acid of formula:\nRxe2x80x2CH2xe2x80x94CO2H\nxe2x80x83with Rxe2x80x2: H or OH,\nxe2x80x83or from 5 to 10% by weight of a mixture of such organic acids;\n(d) from 7 to 13% by weight of anionic surfactant(s), and\n(e) from 3 to 7% by weight of a solubilizing agent consisting of a compound of formula:\nRxe2x80x3O(CH2CH2O)xxe2x80x94H\nxe2x80x83in which\nRxe2x80x3 represents a C2 to C6 alkyl radical,\nx represents an integer from 1 to 3, or a mixture of such compounds.\nThe strong inorganic acid present in the composition is preferably chosen from sulfuric acid, phosphoric acid and their mixtures.\nUse is advantageously made, in the formulation according to the invention, of glycolic acid as organic acid of formula Rxe2x80x2CH2xe2x80x94CO2H or of a mixture of glycolic acid and acetic acid as mixture of organic acids of formula Rxe2x80x2CH2xe2x80x94CO2H. As advantageous example, the mixture of said organic acids contains from 30 to 70% by weight of glycolic acid and from 70 to 30% by weight of acetic acid.\nAccording to a specific embodiment of the invention, the solubilizing agent is chosen from:\nthe compound of formula CH2(CH2)3xe2x80x94Oxe2x80x94CH2CH2OH (2-butoxyethanol or xe2x80x9cbutyl glycolxe2x80x9d)\nthe compound of formula CH2(CH2)3xe2x80x94Oxe2x80x94(CH2CH2O)2H (2-(2-butoxyethoxy)ethanol or xe2x80x9cbutyl diglycolxe2x80x9d).\nIn addition, the composition according to the invention advantageously comprises from 0.1 to 1.0% by weight of a nonionic surface-active agent of general formula:\nRxe2x80x2xe2x80x3O(CH2CH2O)yH\nin which\nRxe2x80x2xe2x80x3 is a saturated C16 to C18 alkyl radical, and\ny is an integer between 8 and 40, or a mixture of such nonionic surface-active agents.\nThis nonionic surface-active agent can consist in particular of a compound of formula:\nRxe2x80x2xe2x80x3O(CH2CH2O)yH\nin which\nRxe2x80x2xe2x80x3 is a saturated C16 to C18 alkyl radical, and\ny is equal to approximately 11.\nSuch a nonionic surface-active agent is sold, for example, by the firm BASF under the name Lutensol AT 11(copyright).\nThe composition advantageously comprises from 0.3 to 0.4% by weight of such a nonionic surface-active agent.\nThe pH of the composition is advantageously less than 5.\nAccording to a specific embodiment, the strong inorganic acid/organic acid(s) of formula Rxe2x80x2CH2xe2x80x94CO2H ratio by weight is between 2/1 and 5/1.\nThe organic acid(s) of formula Rxe2x80x2CH2CO2H/monocarboxylic acid(s) of formula Rxe2x80x94COOH ratio by weight is advantageously between 3/1 and 10/1.\nThe organic acid(s) of formula Rxe2x80x2CH2xe2x80x94CO2H/solubilizing agent of formula Rxe2x80x3Oxe2x80x94(CH2CH2O)xxe2x80x94H ratio by weight is preferably between 3/1 and 1/1.\nUse is advantageously made, as anionic surfactant, of a mixture of anionic surfactants, for example a mixture of a sulfonate surfactant, such as an alkylbenzenesulfonate with 8 to 16 carbon atoms, for example cumenesulfonate, and of a sulfonate from the family of the ether carboxylic acids, for example a mixture of ether carboxylic acids comprising from 15 to 25 carbon atoms and corresponding to the following general formula:\nRxe2x80x3xe2x80x3(OC2H4)nOCH2COOH\nwith Rxe2x80x3xe2x80x3 an alkyl group with 4 to 12 carbon atoms.\nThe sulfonate surfactant is advantageously chosen from alkylsulfonates with 8 to 16 carbon atoms and the mixtures of these. When use is made of a mixture of sulfonate surfactant and of ether carboxylic acid surfactant, use is made of a mixture containing, for example, from 20 to 80% by weight of sulfonate surfactant and from 80 to 20% by weight of ether carboxylic acid surfactant. However, the mixture advantageously comprises more than 50% by weight of sulfonate surfactant.\nAnother subject matter of the invention is the use of a composition according to the invention as fungicidal composition, in particular fungicidal and bactericidal composition.\nA further subject matter of the invention is a process for disinfecting a plant (pipework), vessels, and the like), in which process a composition according to the invention, preferably in dilute form, is brought into contact with the plant.\nThe concentrated composition is advantageously diluted so that the solution, after dilution, has a content of strong inorganic acid of less than 0.4% and at pH of less than 4, preferably of less than 3.\nDisinfecting is advantageously carried out at a temperature of between 0 and 25xc2x0 C."}
{"text": "Current trends in the semiconductor and electronics industry require memory devices to be made smaller, faster and require less power consumption. In addition to active circuitry on such integrated circuit devices, micro-electromechanical (MEMS) devices are sometimes employed in various applications. One exemplary application involves aircraft or vehicle applications, where a MEMS device is employed to detect a change in acceleration. In early systems, the active circuitry was manufactured separately from the MEMS device, however, in recent years attempts have been made to integrate the MEMS device and the active circuitry on the same semiconductor substrate, and thus reside in the same package.\nPrior art MEMS acceleration sensors consist in some instances of a cantilever type beam extending over an underlying layer with a gap therebetween. As the acceleration of a body associated with the sensor changes, the cantilever bends with respect to the underlying layer, causing a change in the gap distance. By employing the cantilever as one element of a capacitor, the change in gap distance results in a change of capacitance, thereby reflecting the change in acceleration.\nPrior art MEMS acceleration sensor systems were not space efficient. That is, the large MEMS portion of the device was fabricated next to the active circuitry employed to detect and communicate the change in capacitance. For example, as illustrated in prior art FIG. 1, an integrated circuit die 10 has a MEMS structure 12 formed next to active circuitry 14, and the MEMS structure 12 occupies a modest amount of die area. This lateral juxtaposition of the MEMS sensor and the active circuitry disadvantageously increases the die size, thereby causing such integrated circuit chips to be rather expensive. Therefore improvements in MEMS sensor devices are desired."}
{"text": "1. Technical Field\nThe present disclosure relates to a server system for conveniently checking, maintaining, or replacing electronic elements.\n2. Description of Related Art\nUsually, a plurality of electronic elements is secured in a server system. So, a server system which is convenient to check, maintain, or replace electronic elements is required."}
{"text": "The present invention relates to the manufacture of substrates, including devices thereon. More particularly, the invention provides a technique including a method and a structure for forming multi-layered substrate structures for the fabrication semiconductor integrated circuits or optoelectronic devices, for example. But it will be recognized that the invention has a wider range of applicability; it can also be applied to other substrates for multi-layered integrated circuit devices, three-dimensional packaging of integrated semiconductor devices, photonic and/or optoelectronic devices (e.g., light valves), piezoelectronic devices, microelectromechanical systems (“MEMS”), nano-technology structures, sensors, actuators, solar cells, flat panel displays (e.g., LCD, AMLCD), biological and biomedical devices, and the like.\nFrom the very early days, human beings have been building useful articles, tools, or devices using less useful materials for numerous years. In some cases, articles are assembled by way of smaller elements or building blocks. Alternatively, less useful articles are separated into smaller pieces to improve their utility. A common example of these articles to be separated include substrate structures, such as a glass plate, a diamond, a semiconductor substrate, and others.\nThese substrate structures are often cleaved or separated using a variety of techniques. In some cases, the substrates can be separated using a saw operation. The saw operation generally relies upon a rotating blade or tool, which cuts through the substrate material to separate the substrate material into two pieces. This technique, however, is often extremely “rough” and cannot generally be used for providing precision separations in the substrate for the manufacture of fine tools and assemblies. Additionally, the saw operation often has difficulty separating or cutting extremely hard and or brittle materials, such as diamond or glass.\nAccordingly, techniques have been developed to separate these hard and or brittle materials using cleaving approaches. In diamond cutting, for example, an intense directional thermal mechanical impulse is directed preferentially along a crystallographic plane of a diamond material. This thermal mechanical impulse generally causes a cleave front to propagate along major crystallographic planes, where cleaving occurs when an energy level from the thermal mechanical impulse exceeds the fracture energy level along the chosen crystallographic plane.\nIn glass cutting, a scribe line using a tool is often impressed in a preferred direction on the glass material, which is generally amorphous in character. The scribe line causes a higher stress area surrounding the amorphous glass material. Mechanical force is placed on each side of the scribe line, which increases stress along the scribe line until the glass material fractures, preferably along the scribe line. This fracture completes the cleaving process of the glass, which can be used in a variety of applications, including households.\nAlthough the techniques described above are satisfactory, for the most part, as applied to cutting diamonds or household glass, they have severe limitations in the fabrication of small complex structures or precision work pieces. For instance, the above techniques are often “rough” and cannot be used with great precision in fabrication of small and delicate machine tools, electronic devices, or the like. Additionally, the above techniques may be useful for separating one large plane of glass from another, but are often ineffective for splitting off, shaving, or stripping a thin film of material from a larger substrate. Furthermore, the above techniques may often cause more than one cleave front, which join along slightly different planes, which is highly undesirable for precision cutting applications.\nAccordingly, certain techniques have been developed to cleave a thin film of crystalline material from a larger donor substrate portion. These techniques are commonly known as “layer transfer” processes. Such layer transfer processes have been useful in the manufacture of specialized substrate structures, such as silicon on insulator. As merely an example, a pioneering technique was developed by Francois J. Henley and Nathan Chung to cleave films of materials. Such technique has been described in U.S. Pat. No. 6,013,563 titled Controlled Cleaving Process, commonly assigned to Silicon Genesis Corporation of San Jose, Calif., and hereby incorporated by reference for all purposes. Although such technique has been successful, there is still a desire for improved ways of manufacturing multilayered structures.\nFrom the above, it is seen that a technique for separating a thin film of material from a substrate which is cost effective and efficient is desirable."}
{"text": "1. Field of the Invention\nThe invention relates to a safety device protecting a nuclear reactor pressure vessel against overpressure failure in the event of inadequate cooling of the core, including a pressure relief valve being pressure-tightly inserted into a wall of the pressure vessel or into a line passing through the wall, the pressure relief valve having a closure element normally blocking a flow from the interior of the pressure vessel to a pressure relief line and being held by a valve closing spring supported against an abutment.\nIn a device for temperature-dependent pressure relief of pressure vessels, which is known from German Patent DE 36 17 524 C2, corresponding to Published European Application No. 0 247 518 A2 and U.S. Pat. No. 4,836,443, a shut-off device provided in a pressure vessel blow-off line is opened at a preset temperature. The effect achieved through this temperature-dependent pressure relief is that in the event of an extremely unlikely failure of all of the cooling devices of a nuclear power station, the pressure in the primary circuit is reduced to a low level before the wall of the reactor pressure vessel can be damaged by an excessive temperature. In the known device according to the applications and patent mentioned above, the temperature sensor is disposed outside the reactor pressure vessel. The known device therefore works with a relatively high time constant, as calculated from the beginning of the impermissibly high temperature to the commencement of the pressure relief response.\nGerman Published, Non-Prosecuted Application DE 35 26 377 A1, corresponding to U.S. Pat. No. 4,777,013, describes a high-temperature reactor having a reactor pressure vessel and a safety valve in the form of a spring valve for limiting the pressure in the reactor pressure vessel in the event of core heating incidents. The reactor pressure vessel is lined internally with a liner connected to a liner cooling system. The valve spring of the safety valve is formed of a material having a spring force which decreases with rising temperature. When the safety valve is open the valve spring is acted on by a flow of outgoing gas and for cooling purposes it is connected to the liner cooling system."}
{"text": "Cloud computing is a way for users to store data and operate computational processes on infrastructure connected by a network. Thus, instead of having to purchase physical infrastructure, users send processes and data out to be run and stored on infrastructure owned by other entities. The user only pays for the amount of data storage or processing capability that he or she desires. This allows the user to tap computing resources that would be impossible without owning actual, physical and vast computing resources.\nBy its very nature, cloud computing gives developers large tool sets for executing programs and tasks along with database storage but it also creates logistical hurdles to implementing those numerous programs, tasks and storage. For instance, developers may be required to write repetitive daemons for cloud managers and object managers to program the logistics each time an object or process is run. This equates to extra work for developers who must devote great resources on cloud logistics, separate and apart from the development of the actual end product."}
{"text": "IPhone Operating System (iOS) device applies iOS as its operating system, which includes Apple products, such as iPhone, iPad, iPod Touch and Apple TV, etc. With rapidly and widely popularizing of iOS devices, the requirements and use of iOS device becomes more and more.\nAfter softwares are distributed, generally the softwares are required to be modified and updated constantly. After being modified or added with a new function by a software developer, the application is released in the form of patches. A user updates those patches, i.e. modifying and updating is accomplished. Software modifying or updating is for satisfying requirements of users and avoiding virus attacking better."}
{"text": "1. Field of the Invention\nThis invention relates to motorcycles, and particularly to three-wheeled motorcycles.\n2. General Description of the Prior Art\nWhile the greater economy achievable by the use of the motorcycle as a basic form of transportation in place of the automobile is widely recognized, there has been a general lack of acceptance of it for this role because of the hazards associated with its operation. Most motorcycles have two wheels, and this causes a much less stable engagement with the road than achievable with a vehicle with three or four wheels. While motorcycles have been constructed with three wheels, typically by adding a rear mounted side car to which is attached a third wheel, this arrangement has not been too satisfactory, particularly because of the asymmetric load and handling characteristics. In recognition of this problem, the applicant previously invented and obtained German Pat. No. 1,063,473 on an improved symmetrically configured three-wheeled motorcycle in which the frame of the vehicle could be banked in unison with all of the wheels. However, there still remained a problem of stability and a sense of insecurity which the applicant believes has slowed the acceptance of three-wheeled motorcycles as safe vehicles. A recently introduced three-wheeled vehicle employs operator controlled actuators for banking. The addition of such controls adds complexity and, of course, costs to the vehicle. More importantly, however, there is a time lag between the operation of the controls and a response by the actuators, and this introduces new hazards if banking commands are not precisely timed.\nAccordingly, it is an object of this invention to provide an improved three-wheeled motorcycle without the need for powered actuators or additional vehicle controls which is a more stable vehicle with demonstrably enhanced safety of operation.\nIt is a further object of this invention to provide a three-wheeled vehicle which is both safe and comfortable in operation in all types of weather."}
{"text": "The processing of paper often involves the crushing of raw materials and the preparation of paper pulp. As there is a lot of manual participation in the process of paper processing, the cost of paper production enterprises remains high, and this device effectively solves this problem."}
{"text": "1. Field of the Invention\nThe invention relates to a mixer oscillator stage. The invention further relates to a mixer oscillator integrated circuit and to a tuner. The invention also relates to a receiver, and more particularly, but not exclusively, a television receiver.\n2. Description of the Related Art\nSuch a mixer oscillator stage is used in, for example, television receivers. Different solutions are known for such a mixer oscillator stage. As is known, television signals comprise signals in the bands historically referred to as UHF and VHF bands. In most cases, the VHF band is split into at least two ranges (VHFL and VHFH).\nFor 3-band tuners, the different bands are often referred to as A, B and C, or Low, Mid and High. For 2-band tuners, it is customary to refer to the bands as VHF-Low, VHF-High and UHF, although the UHF band may be capable of receiving frequencies transmitted below the UHF band.\nThe most straightforward solution is to use the 3-band concept with a separate mixer and oscillator for each band.\nThe above solution is quite complex. To reduce the costs of three different bands, a switching between the low and high VHF signals is performed. The range switching is normally performed with a switching diode which selects a different coil inductance for VHF-Low and VHF-High. Especially, if the mixer oscillator stage is implemented as an integrated circuit which also has a PLL function, referred to as MO-PLL Combi IC, a 3-band tuner is not possible without exceeding the 28-pin structure that the 2-band switched version has. The 28-pin package is one of the largest cost-effective packages available today.\nTo reduce the number of necessary pins, it is proposed to pursue asymmetrical input/output functions. To achieve a good performance, a good isolation between the mixer/oscillator functions, like oscillators, mixers, and the IF amplifier, is necessary. It is hardly possible to achieve this good isolation by using asymmetrical input /output functions.\nBalanced signal pairs and balanced circuitry have so far provided the best performance. The performance must meet the more stringent requirements recently imposed by the European CENELEC directive EN55020. EN55020 is a European directive on EMC requirements for television and sound broadcast receivers and associated equipment. Special attention must be paid during tuner design to ensure that the requirements are met.\nOne of the disadvantages of the known mixer oscillator stages (topologies) is that different mixer oscillator topologies are necessary for the 3-band concept and for the switched 2-band concept, making the development of at least two different topologies expensive and time-consuming. A further disadvantage is that the 3-band concept (especially MOPLL Combi) requires more input and oscillator pins necessary for the mixer oscillator integrated circuit.\nIt is object of the invention to overcome the disadvantages of the prior art and further to provide a mixer oscillator stage, a mixer oscillator integrated circuit, a tuner and a television receiver having an improved performance.\nTo this end, a first aspect of the invention provides a mixer oscillator stage comprising mixing means with at least one mixer, oscillator means with at least two oscillators and switching means coupled between the at least one mixer and the at least two oscillators, the number of oscillators exceeding the number of mixers\nA second aspect of the invention provides a mixer oscillator integrated circuit incorporating the above mixer oscillator stage\nA third aspect of the invention provides a tuner comprising the above mixer oscillator stage.\nA fourth aspect of the invention provides a receiver comprising such a tuner.\nBy choosing number of oscillators to be larger than the number of mixers and by using switching means in between, the same mixer oscillator stage can be used for different applications. The cost of development are thereby largely decreased.\nAn embodiment of a mixer oscillator stage according to the invention has the features that the oscillator means comprises a first, a second and a third oscillator, the mixing means comprises a first and a second mixer, and the switching means are coupled between the mixers and the oscillators.\nBy using three oscillators and switching means, it is possible to use the same mixer oscillator stage for a 2-band switched concept and for a genuine 3-band concept."}
{"text": "Circuitry in state-of-the-art semiconductor chips may be organized into selected regions, blocks or modules containing transistors and other circuit components, such as diodes, resistors, and capacitors, for example. The modules may be separated by channels through which electric wires pass between electric terminations or pins in the modules. The electric wires may be routed through the channels and over or under the modules themselves in separate layers, depending upon the design rules followed by the particular semiconductor chip. Intermediate levels of organization in the modules may include the use of standard cells, for example. In such standard cells, the design and its components are specifically predetermined, and, thus, the circuit characteristics, such as delay times for propagation of a signal through the particular standard cell are well known. Transistors and circuit components may also be organized into non-standard or custom cells which can be referred to as \"leaf cells.\" This terminology is used in the semiconductor industry by analogy to how a tree grows from its trunk into branches and ultimately into leaves, at which photosynthesis takes place. Similarly, the basic circuit chip includes circuitry (the tree) which includes modules (the branches) which are organized into leaf cells (the leaves) which perform circuit activities. Simply stated, the leaf cells are at the lowest level of circuit organization above the transistors and other basic circuit elements. The standard cells themselves can be considered to be leaf cells as well.\nThe placement of leaf cells within electric circuit modules to be fabricated in semiconductor chips is dependent upon the capacitance between leaf cells connected by electric circuit wires and upon the actual worst case delay of signal propagation through particular leaf cells. The capacitance between particular pins on leaf cells to be connected is a function of the length of the connecting wire as well as its per-unit capacitance. Such capacitance can be determined in conventional manner, but there is wide variation between different kinds of leaf cells in terms of actual signal delays through the particular custom circuitry designed into the leaf cell.\nThe determination of leaf cell delays supports transistor level timing-driven place and route operations pursued in the course of developing a semiconductor chip to be fabricated. In standard cells and gate arrays, the delays between inputs and outputs of a particular group of transistors is pre-characterized and delay expressions are readily available, permitting completion of design operations leading to chip manufacture without extensive engineering efforts or without suffering considerable design errors.\nExisting approaches for transistor level place and route operations depend upon representation of the electric circuitry for delay determination with a simple resistor and capacitor network, for example. Such an approach is not sensitive to input and state conditions and patterns occurring in the actual electric circuitry in particular leaf cells."}
{"text": "1. Field of the Invention\nThis invention relates to floor fans and, more particularly to a novel floor fan handtruck apparatus and method for incorporating a floor fan into a handtruck to facilitate transporting one or more floor fans and associated equipment from place to place.\n2. The Prior Art\nFloor fans are well known throughout the janitorial profession. The conventional floor fan is configured around an electrically driven, squirrel-cage blower. The electric motor for the blower is mounted in the interior space of the blower. Screen-enclosed air intakes are placed at each end of the axis of the blower. A discharge chute is located adjacent the bottom of the blower housing and extends outwardly tangentially from the blower housing and parallel to the floor. The discharge chute is designed to allow the operator to direct the blast of air from the floor fan horizontally across the designated portion of the floor. Adjustable risers at the outer end of the discharge chute allow the operator to adjust the angle of the air blast from the discharge chute relative to the floor surface.\nConventionally, a carrying handle is formed in the upper surface of the housing of the floor fan to enable the user to transport the floor fan from place to place. However, experience has shown that a floor fan is awkward and difficult to hand carry particularly for members of the janitorial profession who are less robust or even physically impaired. The result is that there are numerous instances where floor fans are not utilized even though their use would dry the floor faster resulting in the return of the floor to a safer condition more quickly. Further, I have found that the overall bulk of the floor fan causes it to be bumped by the knees of the person carrying it and that the inlet screens have been pushed into contact with the squirrel-cage blower. In view of the foregoing, it would be a significant advancement in the art to provide a cart or handtruck apparatus and method for transporting a floor fan. It would be an even further advancement in the art to incorporate the floor fan into a floor fan handtruck, the floor fan handtruck also serving as a carrier for miscellaneous accessory equipment. It would also be an advancement in the art to provide a floor fan handtruck system for a floor fan whereby a second floor fan handtruck can be releasably attached to the first floor fan handtruck to allow the operator to pull the second floor fan handtruck while pulling the first floor fan handtruck. Such a novel apparatus and method is disclosed and claimed herein."}
{"text": "The original maize plant was indigenous to the Western Hemisphere. The plants were weedlike and only through the efforts of early breeders were cultivated crop species developed. The crop cultivated by early breeders, like the crop today, could be wind pollinated. The physical traits of maize are such that wind pollination results in self-pollination or cross-pollination between plants. Each plant has a separate male and female flower that contributes to pollination, the tassel and ear, respectively. Natural pollination occurs when wind transfers pollen from tassel to the silks on the corn ears. This type of pollination has contributed to the wide variation of maize varieties present in the Western Hemisphere.\nThe development of a planned breeding program for maize only occurred in the last century. A large part of the development of the maize product into a profitable agricultural crop was due to the work done by land grant colleges. Originally, maize was an open pollinated variety having heterogeneous genotypes. The maize farmer selected uniform ears from the yield of these genotypes and preserved them for planting the next season. The result was a field of maize plants that were segregating for a variety of traits. This type of maize selection led to, at most, incremental increases in seed yield.\nLarge increases in seed yield were due to the work done by land grant colleges that resulted in the development of numerous hybrid corn varieties in planned breeding programs. Hybrids were developed by selecting corn lines and selfing these lines for several generations to develop homozygous pure inbred lines. One selected inbred line was crossed with another selected inbred line to produce hybrid progeny (F1). The resulting hybrids, due to heterosis, are robust and vigorous plants. Inbreds on the other hand are mostly homozygous. This homozygosity renders the inbred lines less vigorous. Inbred seed can be difficult to produce since the inbreeding process in corn lines decreases the vigor. However, when two inbred lines are crossed, the hybrid plant evidences greatly increased vigor and seed yield compared to open pollinated, segregating maize plants. An important consequence of the homozygosity and the homogenity of the inbred maize lines is that all hybrid seed produced from any cross of two such elite lines will be the same hybrid seed and make the same hybrid plant. Thus the use of inbreds makes hybrid seed which can be reproduced readily. The hybrid plant in contrast does not produce hybrid seed that is readily reproducible. The seed on a hybrid plant is segregating for traits.\nThe ultimate objective of the commercial maize seed companies is to produce high yielding, agronomically sound plants that perform well in certain regions or areas of the Corn Belt. To produce these types of hybrids, the companies must develop inbreds, which carry needed traits into the hybrid combination. Hybrids are not often uniformly adapted for the entire Corn Belt, but most often are specifically adapted for regions of the Corn Belt. Northern regions of the Corn Belt require shorter season hybrids than do southern regions of the Corn Belt. Hybrids that grow well in Colorado and Nebraska soils may not flourish in richer Illinois and Iowa soil. Thus, a variety of major agronomic traits are important in hybrid combination for the various Corn Belt regions, and have an impact on hybrid performance.\nInbred line development and hybrid testing have been emphasized in the past half-century in commercial maize production as a means to increase hybrid performance. Inbred development is usually done by pedigree selection. Pedigree selection can be selection in an F2 population produced from a planned cross of two genotypes (often elite inbred lines), or selection of progeny of synthetic varieties, open pollinated, composite, or backcrossed populations. This type of selection is effective for highly inheritable traits, but other traits, for example, yield requires replicated test crosses at a variety of stages for accurate selection.\nMaize breeders select for a variety of traits in inbreds that impact hybrid performance along with selecting for acceptable parental traits. Such traits include: yield potential in hybrid combination; dry down; maturity; grain moisture at harvest; greensnap; resistance to root lodging; resistance to stalk lodging; grain quality; disease and insect resistance; ear and plant height. Additionally, Hybrid performance will differ in different soil types such as low levels of organic matter, clay, sand, black, high pH, low pH; or in different environments such as wet environments, drought environments, and no tillage conditions. These traits appear to be governed by a complex genetic system that makes selection and breeding of an inbred line extremely difficult. Even if an inbred in hybrid combination has excellent yield (a desired characteristic), it may not be useful because it fails to have acceptable parental traits such as seed yield, seed size, pollen production, good silks, plant height, etc.\nTo illustrate the difficulty of breeding and developing inbred lines, the following example is given. Two inbreds compared for similarity of 29 traits differed significantly for 18 traits between the two lines. If 18 simply inherited single gene traits were polymorphic with gene frequencies of 0.5 in the parental lines, and assuming independent segregation (as would essentially be the case if each trait resided on a different chromosome arm), then the specific combination of these traits as embodied in an inbred would only be expected to become fixed at a rate of one in 262,144 possible homozygous genetic combinations. Selection of the specific inbred combination is also influenced by the specific selection environment on many of these 18 traits which makes the probability of obtaining this one inbred even more remote. In addition, most traits in the corn genome are regrettably not single dominant genes but are multi-genetic with additive gene action not dominant gene action. Thus, the general procedure of producing a non segregating F1 generation and self pollinating to produce a F2 generation that segregates for traits and selecting progeny with the visual traits desired does not easily lead to a useful inbred. Great care and breeder expertise must be used in selection of breeding material to continue to increase yield and the agronomics of inbreds and resultant commercial hybrids.\nCertain regions of the Corn Belt have specific difficulties that other regions may not have. Thus the hybrids developed from the inbreds have to have traits that overcome or at least minimize these regional growing problems. Examples of these problems include in the eastern corn belt Gray Leaf Spot, in the north cool temperatures during seedling emergence, in the Nebraska region CLN (corn Lethal necrosis and in the west soil that has excessively high pH levels. The industry often targets inbreds that address these issues specifically forming niche products. However, the aim of most large seed producers is to provide a number of traits to each inbred so that the corresponding hybrid can useful in a broader regions of the Corn Belt. The new biotechnology techniques such as Microsatellites, RFLPs, RAPDs and the like have provided breeders with additional tools to accomplish these goals."}
{"text": "1. The Field of the Invention\nThe present invention relates to vacuum shrouds. More specifically, the present invention relates to vacuum shrouds for angle grinders which include an access hatch and a retention mechanism for the access hatch.\n2. State of the Art\nDust shrouds are typically used to both contain grinding dust and prevent a mess as well as to protect an operator from debris, dust, and the grinding disc. These shrouds generally attach to an angle grinder between the grinder body and the grinding wheel. The shroud then forms a barrier between the operator and the grinding wheel. The shroud often extends around the edges of the grinding wheel, leaving the bottom grinding portion open for use. Thus, the dust shroud provides a dust and debris barrier and a safety barrier.\nSome dust shrouds provide a vacuum attachment so that dust and debris can be suctioned out of the shroud. This avoids airborne dust and debris that can lead to such health problems as silicosis. The use of a vacuum also keeps the workplace clean of debris, dramatically reducing the time necessary to clean up after a grinding job. Vacuum dust shrouds are commonly used while grinding concrete. It can be appreciated how, without a shroud, a significant amount of concrete dust is spread across a wide area. With a vacuum shroud, very little dust goes uncollected. The use of a vacuum shroud can also protect the tool by preventing dust and debris from entering the motor and gears, thus extending tool life.\nEffective vacuum shrouds often cover the top and entire perimeter of the grinding disc, such that most of the dust and debris is captured by the vacuum and the disc is not exposed to the operator. Complete enclosure of the perimeter of the grinding disc makes it difficult to use against a corner, such as when grinding a floor adjacent a wall.\nVacuum shrouds for grinders have been made which allow a user to remove a hatch portion of the shroud to expose a side of the grinding disk and allow grinding into a corner. It has been found, however, that workers tend to remove the hatch to grind a corner and either lose the hatch or neglect to replace the hatch after finishing the corner. The worker needs both hands for operating the grinder, and thus can not hold the hatch while grinding. It is not uncommon for workers to simply place the hatch in a toolbox or on the floor, and lose the hatch or simply not replace the hatch as is desired.\nThe failure to replace the hatch when not grinding a corner is more detrimental than one might think. In addition to the danger of having an exposed grinding disk, having an opening in the perimeter of the dust shroud allows a significant portion of the grinding dust to escape and not be collected by the vacuum. The dust from grinding concrete is very fine, and will float in the air for a long period of time, allowing uncollected dust to settle over a large area. Additionally, the dust contains compounds from the concrete which are harmful to the lungs, and uncollected dust poses more of a health risk to the worker and to others.\nThus there is a need for a grinder dust shroud with a removable corner grinding hatch which encourages workers to use the hatch. There is a need for a grinder shroud which allows a person to store an access hatch on the shroud itself when the hatch is removed from the shroud."}
{"text": "Heretofore, it has been suggested to use a learning device consisting of a combination of a learning sheet, learning reactor, and answer indicating bar. A base plate is made by arranging one common conducting band and separate single conducting bands set through slight clearances respectively to the right and left of it linearly in a series over a plurality of streaks by a printed wiring technique on an insulating plate. A right answer insulating part is formed between the common conducting band and one single conducting band. A wrong answer insulating part is formed between the common conducting band and the other single conducting band, so that respectively different indications may be made in the respective insulating parts. A plurality of answer holes are formed on the learning sheet mounted on the base plate so that only one answer hole representing a right answer is present in the right answer insulating part, and a plurality of answer holes representing wrong answers are present in the wrong answer insulating part. The answer hole considered to be the right answer to the set question is selected from among the answer holes. The conducting bands present in this answer hole are short-circuited by contacting the conductor of the answer indicating bar with this selected answer hole. An electric indication, i.e., the lighting of an indicating lamp, as to whether the selected answer is right or wrong is made. In the case of the right answer, a blue lamp is continuously lighted. In the case of the wrong answer, a red lamp is intermittently lighted. Or a buzzer is sounded continuously in the case of the right answer, and intermittently in the case of the wrong answer.\nHowever, in this suggested learning device, the respective conducting bands arranged on the base plate of the learning reactor or the respective insulating parts of the right answer and wrong answers are formed in straight lines and are exposed out of the answer holes on the learning sheet. Therefore, the learner is likely to memorize the arranged patterns of the conducting bands. Because the respective answer holes to two or more set questions on one straight line all show the same reaction indications, even if the set question is not understood, the right answer may be selected and this negates objective learning.\nBecause the answer hole on the learning sheet is selected and the respective hole on the conducting bands are short-circuited by the contact of the conducting parts through this answer hole, the learning sheet itself is not likely to be damaged and there is an advantage in that the learning sheet can be repeatedly used. However, because the answer is obtained by the short-circuiting input between the conducting bands by the conductor of an answer indicating bar, it is necessary to use a conductive metal member for the conductor and to positively contact the conducting bands. To short-circuit the conducting bands adjacent to each other, this conductor must be contacted at a specific angle. When this angle is inclined even slightly, the contact will be faulty and no answer input will be attained. Particularly, as such learning device is to be used mostly by young school age children, the operation will be rather rough, the defects will be conspicuous, no proper answer indication will be obtained, and the right and wrong answers will be indicated by mistake.\nFurther, in order to use this kind of learning device, e.g., for learning a foreign language, this learning device includes a sight teaching material connected together with a tape recorder as a hearing teaching material. Heretofore, the sight learning with the learning reactor and learning sheet, and the hearing learning with the recorded sound reproduction of the tape recorder, have been used separately. Therefore, the thinking time on the part of the learner, and the timing of the set question and others, are all determined directly by the driving manner of the tape. Not only is this an inconvenience in that the individual difference of each learner cannot be inserted into the operation, but also, in the case of a set question requiring thought time, it will be necessary to stop the drive of the tape each time by a manual operation, or to adopt such means as of leaving a blank corresponding to the thinking time on the tape after setting the question. Stopping the drive of the tape each time is a useless complication to the learner, and will prevent the learner from concentrating solely on the learning. Providing blanks on the tape results in tape loss. Further, when replacing the learning sheet, the tape drive controlling operation will be forced on the learner and, as a result, will be likely to obstruct the learning."}
{"text": "The invention relates to an exercise and therapeutic apparatus and more specifically a rebound system and an improved method of altering the resistance of the carriage that the exercise participant travels upon.\nThe Invention relates to the art of exercise apparatuses that generally consists of a stationary frame, a moving carriage, a kickplate, a resistance system and a rebound system. The exercising orientated is positioned in a supine position on the carriage as it travels in a longitudinally oscillating or reciprocating motion by exerting a force upon a kickplate that is attached to the foot end portion of the frame or by pulling upon arm cords connected head end portion of the frame.\nPrior exercise apparatuses that have been disclosed in U.S. Pat. Nos. 5,042,797 and 5,364,327 that are also applications made by the inventor, disclose resistance systems that allow adjustability from the foot end portion (the portion where the kickplate is located) of the exercise apparatus only. This required the exercising participant or therapist to adjust the elastic members near the kickplate that is located to the foot end portion of the frame.\nMore advanced exercising participants exert greater force upon the kickplate therefor reposition the frame in a longitudinally footward direction. Hence, the acceleration of the mass of the carriage and exercising participant in the longitudinally headward direction creates a counteracting force upon the kickplate and frame in the longitudinally footward direction. This force is great enough to overcome the frictional forces between the lower ground contact portions of the frame and the surface the exercise apparatus rests upon. To help counteract the longitudinally footward movement of the frame, oftentimes the foot end portion of the apparatus was positioned against a wall. This allowed the longitudinally footward force exerted upon the kickplate and frame to be transferred to the rigid wall.\nThis presents a new problem for the individual changing the resistance exerted upon the carriage. To increase the resistance additional elastic members are grasped from the foot end portion of the carriage and then attached to the foot end portion of the frame. Likewise, to decrease the resistance elastic members that are attached to the extreme foot end portion of the frame are disengaged from the frame and the head or handle portions of the elastic members are returned to the foot end portion of the carriage. Because the foot end portion of the machine was usually positioned against a wall, to add additional elastic members the exercising participant or therapist had to reach under the kickplate and grab the head portions of the elastic members from the carriage and pull in a longitudinally footward direction to a mounting bracket on the extreme foot end portion of the frame. This was an awkward movement that was difficult to execute and other times was not feasible because of the immediate location of the wall.\nAnother issue with the earlier designs was that the elastic members that were not in an operative position traveled with the carriage. This added extra mass to the carriage which in turn created greater forces upon the exercising participant during accelerations.\nA further issue with the earlier designs was that the elastic members would develop cracks on the interior surfaces near the head or handle portions. The structure of the handle portion of the elastic members consisted of a handle member that has a barbed longitudinally extending member with a smaller diameter that is frictionally engaged within the end portion of the hollow rubber tube. When the person adjusting the elastic cord grabbed on to the handle portion of an elastic strap oftentimes they would not pull the handle portion substantially longitudinal in an aligned manner with the longitudinal axis of the elastic member. But rather, they would apply a moment perpendicular to the longitudinal axis of the elastic member. This caused a slight rotation of the head portion about an orthogonal axis to the longitudinal axis of the elastic member. The effect of this moment was that it caused the barbed insert of the head member to gouge into the interior surface of the elastic member. The effect of the gouging of the interior surface was most pronounced at the deepest inserted portion of the barbed member that was inserted into the elastic member. After many cycles of changing elastic members from an inoperative positioned to an operative position, the elastic members would break as a result from a crack that began on the interior surface near the end portion of the barbed insert of the head member.\nThere is further an improved rebound system over the prior art devices (namely the assembly discussed in U.S. Pat. No. 5,042,797 discussed further below). The prior art rebound systems were effective in storing energy and not allowing excessive accelerations upon the carriage when the exercising participant traveled beyond the intended longitudinal range of motion. However, the prior art rebound systems were difficult to adjust in the longitudinal range of motion of travel of the carriage, and further, they were more expensive to produce.\nRelevant prior art is disclosed in three earlier U.S. patent applications by the present inventor. The earlier applications are U.S. Ser. Nos. 696,254, 786,540, and 001,192. The apparatus of U.S. Ser. No. 001,192 comprises a platform or carriage, termed a shuttle, being mounted on rollers which engage the tracks of a track assembly or frame. The shuttle is further interconnected to the track assembly with primary and secondary energy storage and release systems comprising elastic cords. The tension forces in the cords of both systems are adjustable to influence the at-rest location of the shuttle on the tracks and the forces applied to the shuttle by the cords as the shuttle moves with respect to the tracks.\nThe primary system interconnects between the shuttle and both ends of the track assembly. The secondary energy absorption and release system is in constant engagement with the track assembly and intermittent engagement with the shuttle, the shuttle engaging the secondary system when the shuttle approaches one or both the extreme of its reciprocating motion on the tracks.\nThe engagement between the secondary system, specifically the rebound assembly, and the track assembly is effected by engagement of projections from the rebound assembly with structural stop members attached to the track assembly. The stop members are positioned so that they are engaged by the projections as the shuttle nears the limits of its travel along the tracks. After engagement the rebound assembly stopped while the shuttle continues to move, stretching the elastic cord(s) and producing forces which decelerate, stop and re-accelerate the shuttle in the opposite direction or assist the primary system in doing so if the primary system is set to operate at that point.\nA more significant piece of prior art known to the inventor is disclosed in U.S. Pat. No. 5,364,327 that is issued to him. This application discloses an exercise apparatus having a frame that comprises two longitudinally extending tracks and a kickplate. The apparatus further has a carriage which is designed to move longitudinally in a reciprocating motion on the frame. The primary energy storage and release system comprises elastic cords attached to the bottom of the shuttle and to the ends of the track assembly at various points along the cords, using ferrules on the cords engaging slots in the end structures of the track assembly. These cords may be connected, and thus constantly engaged during operation, between the shuttle and one end of the track assembly, the shuttle and the other end of the track assembly or the shuttle and both ends. The system functions to bias the shuttle and its occupant toward a position along the tracks between its ends, the forces in the cord(s) serving to decelerate, stop and reaccelerate the shuttle at each end of its travel along the tracks, aided by the secondary system when necessary.\nAnother relevant piece of prior art is illustrated in U.S. Pat. No. 6,042,523, issued Mar. 28,2000, entitled xe2x80x9cTherapeutic Exercise Apparatus and Methodxe2x80x9d. This invention shows a therapeutic exercise devises that is adapted to be positioned on a table. The device utilizes elastic members that are mounted to the carriage.\nIt is therefore the object of the invention to provide a head end positioned resistance adjustment system that has a more convenient access to the person adjusting the resistance of the exercise apparatus.\nIt is another object of the invention to remove the need of having the elastic members that are not in an operative position travel with the carriage by permanently attaching the elastic members to the frame portion and attaching the head portion of the elastic members to the carriage when additional exercise resistance is required.\nIt is further an object of the present invention to provide a head end positioned resistance adjustment system that utilizes a flexible strap where a downward force and an upward force upon the strap at a first and second location will reposition a corresponding elastic member to an operative position.\nIt is another object of the present invention to provide a resistance system that is easy to maintain and replace damaged elastic members.\nIt is another object of the invention to provide a rebound system that is easy to manufacture and provides an exponentially increasing force in the longitudinal direction upon the carriage with respect to longitudinal travel of the carriage at the extreme longitudinal locations of the carriage\"\"s range of motion.\nIt is another object of the invention to provide a rebound system that allows the carriage\"\"s range of motion to be adjusted in a convenient manner.\nFurther objects and advantages of the present invention will become apparent within the detailed description of the present invention.\nThe present invention relates to an improved resistance system that is to be implemented in an exercise apparatus which comprises a frame and a carriage. The frame has a rearward or foot end portion, a head end or forward portion, a leftward lateral side and a rightward lateral side. The carriage has a foot end portion, a head end portion, a leftward lateral side and a rightward lateral side and is adapted to travel longitudinally along the frame. An engagement section is located on the carriage. The resistance system utilizes a plurality of elastic members that extend headwardly from a base portion located in the foot end portion of the frame. The elastic members have a head portion which is on the longitudinally opposite side of the base portion and these head members are adapted to attach to the engagement section of the carriage. When the elastic members are attached to the engagement section this is referred to an operative position where tension in the longitudinally footward direction is applied to the carriage.\nThe present invention further has a rebound system that is to be implemented in the exercise apparatus where the rebound system comprises two flexible members. The first flexible member has a first portion, a central portion and a second portion where the first portion is mounted to the carriage and the second portion is mounted to the frame. Likewise, the second flexible member has a third portion, a central portion and a fourth portion where the third portion is mounted to a position on the carriage that is laterally right of the first portion of the first flexible member. The fourth portion is mounted to the frame laterally right of the second portion of the first flexible member. A resistance member is positioned substantially in the central portions of the first and second flexible members. The resistance member resists lateral separation of the central portions of the first and second flexible members that in turn causes tension in the first and second flexible members which exerts an exponentially increasing force upon the carriage when the carriage reaches its extremities in longitudinal travel."}
{"text": "1. Field of the Invention\nThe present invention relates to the processing of audio signals and, more particularly, the coding of multi-channel audio signals.\n2. Description of the Related Art\nParametric multi-channel audio coders generally transmit only one full-bandwidth audio channel combined with a set of parameters that describe the spatial properties of an input signal. For example, FIG. 1 shows the steps performed in an encoder 10 described in International Application No. WO2003/90208, filed Apr. 22, 2003.\nIn an initial step S1, input signals L and R are split into subbands 101, for example, by time-windowing followed by a transform operation. Subsequently, in step S2, the level difference (ILD) of corresponding subband signals is determined; in step S3, the time difference (ITD or IPD) of corresponding subband signals is determined; and in step S4, the amount of similarity or dissimilarity of the waveforms which cannot be accounted for by ILDs or ITDs, is described. In the subsequent steps S5, S6, and S7, the determined parameters are quantized.\nIn step S8, a monaural signal S is generated from the incoming audio signals, and finally, in step S9, a coded signal 102 is generated from the monaural signal and the determined spatial parameters.\nFIG. 2 shows a schematic block diagram of a coding system comprising the encoder 10 and a corresponding decoder 202. The coded signal 102, comprising the sum signal S and spatial parameters P, is communicated to a decoder 202. The signal 102 may be communicated via any suitable communications channel 204. Alternatively, or additionally, the signal may be stored on a removable storage medium 214, which may be transferred from the encoder to the decoder.\nSynthesis (in the decoder 202) is performed by applying the spatial parameters to the sum signal to generate left and right output signals. Hence, the decoder 202 comprises a decoding module 210 which performs the inverse operation of step S9 and extracts the sum signal S and the parameters P from the coded signal 102. The decoder further comprises a synthesis module 211 which recovers the stereo components L and R from the sum (or dominant) signal and the spatial parameters.\nOne of the challenges is to generate the monaural signal S, step S8, in such a way that, on decoding into the output channels, the perceived sound timbre is exactly the same as for the input channels.\nSeveral methods of generating this sum signal have been suggested previously. In general, these methods compose a mono signal as a linear combination of the input signals. Particular techniques include:\n1. Simple summation of the input signals. See, for example, ‘Efficient representation of spatial audio using perceptual parametrization’, by C. Faller and F. Baumgarte, WASPAA′01, Workshop on applications of signal processing on audio and acoustics, New Paltz, New York, 2001.2. Weighted summation of the input signals using principle component analysis (PCA). See, for example, International Patent Application No. WO2003/85645, filed Mar. 20, 2003 and International Patent Application No. WO2003/85643 filed Mar. 20, 2003. In this scheme, the squared weights of the summation sum up to one and the actual values depend on the relative energies in the input signals.3. Weighted summation with weights depending on the time-domain correlation between the input signals. See for example ‘Joint stereo coding of audio signals’, by D. Sinha, European Patent Application No. EP 1 107 232 A2. In this method, the weights sum to +1, while the actual values depend on the cross-correlation of the input channels.4. U.S. Pat. No. 5,701,346 to Herre et al. discloses weighted summation with energy-preservation scaling for downmixing left, right, and center channels of wideband signals. However, this is not performed as a function of frequency.\nThese methods can be applied to the full-bandwidth signal or can be applied on band-filtered signals which all have their own weights for each frequency band. However, all of the methods described have one drawback. If the cross-correlation is frequency-dependent, which is very often the case for stereo recordings, coloration (i.e., a change of the perceived timbre) of the sound of the decoder occurs.\nThis can be explained as follows: For a frequency band that has a cross-correlation of +1, linear summation of two input signals results in a linear addition of the signal amplitudes and squaring the additive signal to determine the resultant energy. (For two in-phase signals of equal amplitude, this results in a doubling of amplitude with a quadrupling of energy.) If the cross-correlation is 0, linear summation results in less than a doubling of the amplitude and a quadrupling of the energy. Furthermore, if the cross-correlation for a certain frequency band amounts −1, the signal components of that frequency band cancel out and no signal remains. Hence, for simple summation, the frequency bands of the sum signal can have an energy (power) between 0 and four times the power of the two input signals, depending on the relative levels and the cross-correlation of the input signals."}
{"text": "1. Field of the Invention\nThe present invention relates to an active energy ray-curable composition and an antistatic film.\n2. Description of the Related Art\nGenerally, resin materials exhibit excellent electrical insulation properties and are therefore very useful for applications requiring electrical insulating properties such as insulators, but surfaces thereof are easily charged with static electricity, which is thus likely to result in dust adsorption and electrostatic troubles.\nTo solve such problems, it is known to use various antistatic agents (in particular, cationic antistatic agents).\nFor example, JP2011-225797A discloses an “antistatic coating composition including an antistatic copolymer (A) having a quaternary ammonium base in the molecule, a compound (B) having three or more ethylenically unsaturated groups, a photopolymerization initiator (C) containing a photopolymerization initiator (c-1) having a solubility in water of 0.2 g/L or less and a photopolymerization initiator (c-2) having a solubility in water of 1 g/L or more, and a solvent (D) containing an alcohol-based solvent (d-1)” ([claim 1]), and also discloses that the photopolymerization initiator (C) is “at least one selected from the group consisting of 1-hydroxy-cyclohexyl-phenyl-ketone, 2-methyl-1-[4-(methylthio)phenyl]-2-morpholinopropan-1-one, bis(2,4,6-trimethylbenzoyl)-phenylphosphine oxide, 2-hydroxy-1-{4-[4-(2-hydroxy-2-methyl-propionyl)-benzyl]-phenyl}-2-methyl-propan-1-one, 2,4,6-trimethylbenzoyl-diphenyl-phosphine oxide, oligo {2-hydroxy-2-methyl-1-[4-(1-methylvinyl)phenyl]propanone}, and 2,4,6-trimethylbenzophenone” ([claim 4])."}
{"text": "Retrofitting non-digital billboards have proven to be expensive, time consuming and labor intensive. Moreover, simply removing an older non-digital billboard and replacing it with a new digital billboard has not proven entirely satisfactory since older installed, non-digital, billboard panels represent substantial capital outlays making it financially difficult, if not impossible, to discard such panels arbitrarily for replacement with digital panels. Therefore, it would be highly desirable to have a new and improved billboard retrofit kit that can be easily and quickly installed on any structural surface, such as an existing billboard, without the need to replace or discard existing media mounting panels. The new and improved billboard retrofit kit should greatly improve displayed information, displaying such advertising information, with improved resolution, contrast and brightness characteristics. Moreover, the retrofit kit should enable the displayed content to be easily and quickly changed or updated, either on-site or remotely, at a lesser cost than updating the content of an older billboard. Finally, installation of the kit in the field on an existing billboard structure should not require any special installation equipment and should be able to be accomplished by one or two individuals in a fast and convenient manner."}
{"text": "In the past, in order to carry a photocatalyst on a carrier composed of an organic material such as resin, a silicone-modified resin, or resin containing a fixed amount of colloidal silica or polysiloxane, was known to be dissolved in an organic solvent such as xylene, and the dissolved resin was first coated onto a carrier to form an adhesive layer (see, for example, Patent Document 1)."}
{"text": "Polymeric thin films containing asymmetrical chromophores have been under investigation for their second-order nonlinear optical (NLO), piezoelectric, and/or pyroelectric properties for over ten years. This area of research has yielded several types of versatile polymers useful for applications such as optical second harmonic generation and electro-optic modulation of optical signals. The first generation organic NLO polymers (NLOPs) were comprised of small asymmetrical chromophores dissolved in glassy polymers (guest-host systems). The materials science has since evolved to sidechain and mainchain systems in which the asymmeteric chromophore is chemically attached to the polymer, as discussed in more detail below.\nOptical nonlinearity\nTo a first approximation, it is the molecular structure of the chromophore and its orientation that dictate the nonlinear optical properties of the system, and it is the polymer structure that dictates the processability and temporal stability of the final product. Second-order nonlinear optical films do not have a center of symmetry. In order to understand the symmetry requirements, one must consider the three-dimensionality of the electric polarizability in the material, the directions of the applied electric fields, and the electric fields of light rays passing through the material. A treatment most often used to describe this interaction is the following nonlinear polarization equation (a power series in electric field strength, E: EQU P=.sub..chi..sup.(1) E+.sub..chi..sup.(2) EE+.sub..chi..sup.(3) EEE+ . . .\nwherein P is the polarization in the material, .chi..sup.(1) is the coefficient for linear interaction with E, .chi..sup.(2) is the second-order nonlinear optical coefficient, and .chi..sup.(3) is the third-order nonlinear optical coefficient. The second-order nonlinear optical coefficient is a third-rank tensor whose tensor elements are all zero in materials possessing the center-of-inversion symmetry.\nIn order for organic films to have large second-order nonlinear optical coefficients, .chi..sup.(2), they must contain a high concentration of asymmetrical, highly hyperpolarizable chromophores, arranged in a highly polarized configuration. NLOP chromophores must have an inherently large molecular second-order nonlinearity, .beta. (also called second-order susceptibility, quadratic susceptibility, and first hyperpolarizability).\nFor predicting the magnitude of the first molecular hyperpolarizability coefficient beta (.beta.), to a first approximation, considering only the ground state of the chromophore and the first excited singlet state leads to a very successful two-state model. The most important requirement for a large beta is that the asymmetrical chromophore has a low energy of transition between ground and first excited singlet state. It is also necessary that there be a large change in dipole moment in going between the ground state to the first excited singlet state, and that the transition dipole moment be large. In general, large II-conjugated links (e.g., vinylenes, phenylenes, thienylenes, etc.) connecting properly balanced electron donor and acceptor groups are desirable. In this disclosure, \"NLOP\" is used only in reference to polymers that can be formed into films that exhibit significant second-order optical nonlinearity, that is, the NLOP .chi..sup.(2) is at least three times greater than the .chi..sup.(2) of quartz which is about 1.0 pm/V.\nThe following mathematical treatment can be applied to an isolated chromophore's dipole moment: .mu.=.mu..sub.o +.alpha.E+.beta.EE+ . . . , where .mu. is the dipole moment upon applying the electric field E, .mu..sub.o is the ground state dipole moment, .beta. is a third rank tensor describing the chromophore's second order susceptibility, which is also called the quadratic hyperpolarizability or first hyperpolarizability.\nGuest-Host NLOP\nThe simplest form of nonlinear optical polymer, called guest-host (G-H) systems, are solid solutions of small chromophoric molecules and high molecular weight polymers. Generally the G-H system contains about 10 to 20% by weight of the chromophore (higher levels tend to phase separate or have lower glass transition temperatures due to plasticization). The chromophores in G-H systems are somewhat labile (they diffuse by translation and rotation) and they evaporate or sublime at elevated temperatures. Small chromophores, which are easily absorbed through the skin, can be toxic, mutagenic, teratogenic, and carcinogenic. High molecular weight polymers cannot be absorbed through the skin. Therefore, by attaching the chromophore to the polymer, the lability and toxicity problems are solved.\nSidechain NLOP\nSidechain (SC) polymers are comprised of assymetric chromophores which are chemically linked or tethered to the backbone at one end of the chromophore, e.g., at the electron accepting end or at the electron donating end, and the majority of the chromophore is pendent to the backbone. Most of the early reported SC polymers were derived from the free radical copolymerization of (meth) acrylic functional chromophores and methyl methacrylate. Some SC polymers have been designed for Langmuir-Blodgett deposition, e.g., the polymers described in U.S. Pat. No. 5,162,453, issued Nov. 10, 1992 to Hall et al, and U.S. Pat. No. 5,225,285, issued Jul. 6, 1993 to the same inventors, which are described below and incorporated herein by reference.\nMainchain NLOP\nMainchain (MC) polymers are comprised of asymmetrical chromophores which are linked in the backbone at both ends of the chromophore, and the majority of the chromophore forms part of the backbone, the rest of the backbone being connecting groups, bridging one chromophore to the next. The asymmetrical chromophores can be linked in a head-to-tail configuration (isoregic), in a head-to-head configuration (syndioregic), or randomly head-to-tail and head-to-head (aregic). The mainchain chromophoric topology is inherently more stable than the sidechain topology because the chromophore has one less degree of freedom of motion.\nStability\nGeometric considerations, such as how to orient the asymmetrical chromophore in a polar alignment, how to increase the glass transition temperature of the NLOP, and how to prevent relaxation of the chromophore in NLOPs are of utmost importance for improving long term optical stability. There are two types of stability. Physical stability refers to stability of alignment of the chromophore. Chemical stability refers to the integrity of the original chemical structure of the chromophore. Chemical changes can be brought about by thermal energy, electromagnetic energy (photons), reactions with oxygen or other chemicals, or any combination of the above. Chemical changes can range from isomerization (cis--trans), to cyclization, to oxidation, to bond cleavage, or to free radical additions, all of which change the refractive index and .chi..sup.(2) properties. The temporal stability of a .chi..sup.(2) NLOP polymer film refers to how well the physical stability(alignment) and chemical stability of the chromophores are maintained at a given temperature. The processing stability refers to how well the polymer holds up under film processing and packaging conditions.\nIn order to mass produce low cost integrated devices, reliable thin film processing techniques are necessary for integration with silicon chip manufacturing. If these polymers can be easily processed, their uniquely low dielectric constant and large second-order nonlinear optical coefficient (.chi..sup.(2)) will ensure their strong competition for high speed optical signal processing applications. Several processes for giving polar order to the chromophores in a polymer film exist. These include electric-field poling and self-assembly techniques, such as Langmuir-Blodgett (LB) processing.\nElectric-Field Poling\nThin films of the polymer are prepared for poling by spin-coating a liquid solution containing about 10% of the polymer onto a solid substrate. The solvent is removed by baking the film just above the glass transition temperature (T.sub.g). An electric field is applied across the film by corona poling the film as it sits on a grounded conductor plane, or by charging two electrodes contacting the film. There may also be thin cladding or buffer layers between the electrodes and the nonlinear optical polymer. The films are typically poled at a field strength of about 50 to several hundred volts per micron. The electric field is applied at or above T.sub.g for 10 to 60 minutes, then the film is cooled with the field on. After the external field is removed, a net alignment of dipole moments can remain essentially locked in the film for years as long as the temperature of the film remains well below any solid state transition, such as the T.sub.g. This imparts noncentrosymmetry to the film which is necessary for its second-order nonlinear optical properties.\nThere are major problems with electric-field poling. The polymer must be heated to quite high temperatures at which the Brownian motion and rotation of the molecules works against the torque exerted by the electric field on the dipole moment of the chromophores. Hence, at higher temperatures, the chromophores will be less well ordered (i.e., the degree of order .apprxeq..mu.E/kT, where .mu. is the ground state dipole moment of the chromophore, E is the applied electric field, k is Boltzmann's constant, and T is the temperature). Polymers containing formal charges are difficult to pole with an electric field because the charges tend to migrate through the polymer and short-out the electrode. It is difficult to have layers with opposed orientations of the chromophores within the same polymer film as is sometimes desirable for phase matching in frequency doubling applications.\nWhen chromophores in a molten polymer film are exposed to an electric field, the torque of the applied field on the dipoles of the chromophores aligns the average direction of their ground state dipole moments parallel to the direction of the applied field. The largest component of the .beta. tensor of a chromophore is usually at some angle away from the direction of the ground state dipole moment of the chromophore, but \"goes along for the ride\" and becomes aligned (on the average) as the dipoles are aligning with the electric field. .chi..sup.(2) is a tensor quantity describing the second-order nonlinear optical properties of the film, and is a result of the sum of the .beta. tensors of the individual chromophores. Hence, for electric-field-poled films, the .chi..sup.(2) tensor has conical symmetry around the direction of the applied electric field. Thus, the largest component of the .chi..sup.(2) tensor (e.g., d.sub.33 for corona poling) results from those parts of the largest components of the .beta. tensors that are projected onto the direction of the applied electric field. Estimates of .beta. from molecular orbital calculations, combined with estimates of the conical angle about the poling direction from birefringence measurements, have resulted in reasonably accurate d.sub.33 predictions.\nLangmuir-Blodgett (LB) Processing\nTo make films by LB processing, an organic compound is floated on a liquid, usually water, in a trough which holds the water. A solid substrate is dipped through the air-water interface depositing on the substrate a single molecular layer. Thicker films comprised of multilayers of polymer are built up by dipping the substrate repeatedly into and/or out of the LB trough, depositing one molecular layer per stroke. Turn-key, computer-automated, multi-compartment troughs are available from many commercial suppliers, such as NIMA (Coventry CV4 7EZ, England), NLE (Nagoya, 468, Japan), and KSV (SF-00380 Helsinki, Finland).\nIn LB processing, the polymer molecules are designed to have hydrophilic and hydrophobic groups which cause the polymer to float in one conformation on the water. These hydrophilic/hydrophobic forces are useful for removing the centrosymmetry by orienting the asymmetrical chromophores with respect to the plane of the film.\nMultilayer LB films can be formed in three different configurations. Historically these are called \"X\"-, \"Y\"- and \"Z\"-type films, where X is made by depositing always on the down-stroke, Z is made by depositing always on the up-stroke, and Y is made by alternating up- and down-strokes. For the case in which the large dipole moment of the chromophore is normal to the plane of the substrate, all-up or all-down films will be polarized, and the up-down films will not be polarized (dipoles in adjacent layers cancel out). However, in Y-type deposition, sometimes the chromophores will align in a herring-bone pattern due to shear forces in the dipping direction caused by pulling and pushing the substrate through the film floating on the water or by local packing considerations. The herring-bone pattern results in a net polarization in the plane of the film, and can be detected by the generation of second harmonic light. The Y configuration is thermodynamically more stable (e.g., sometimes X and Z configurations spontaneously rearrange in the solid state to the Y configuration). Another way to make polar films by Y-type deposition is to alternately interleave asymmetrical chromophoric polymer layers with optically inert polymer layers. Or one may use a two-compartment LB trough, building up bilayers in which the two polymers are configured with the charge transfer axis of their chromophores respectively pointing in the opposite sense with regard to the hydrophilic and hydrophobic parts of the polymer to which they are attached.\nMolecular weight (and its distribution) is one of several key parameters which affect the physical behavior of polymer molecules in monolayer films at the air-water interface (Langmuir films). The modulus and viscosity of Langmuir films under surface compression vary greatly with molecular weight.\nLB processing has the advantage over electric-field poling in that it can be done at room temperature (or lower); hence, the kT Brownian motion is much less. Furthermore, formal ionic charges on the polymer will not hinder the ordering process; in fact, they can be taken advantage of in designing hydrophilicity into the polymer. This film processing procedure has been reduced to practice for sidechain polymers, but not for mainchain chromophoric polymers. However, mainchain polymers are inherently more stable than sidechain polymers because both ends of the chromophore are attached to the backbone of the polymer, whereas only one end of the chromophore is attached to the backbone in a sidechain polymer, and the other end is dangling free at one of the interfaces of each monolayer. Furthermore, only mainchain polymers that contain the asymmetrical chromophore in a head-to-head (syndioregic) configuration can have the chromophore disposed in a polar alignment normal to the plane of the film by hydrophilic/hydrophobic forces. This is because a mainchain polymer containing chromophores in a head-to-tail configuration will have its chromophores lying on the water, in the plane of the film, and the chains can be pointing in opposite directions. The utility of the syndoregic configuration for L.B. processing will become clear in the detailed description of the invention below.\nU.S. Pat. No. 5,162,453, issued Nov. 10, 1992 to Hall et al, and U.S. Pat. No. 5,225,285, issued Jul. 6, 1993 to the same inventors, describe a multilayered polarized film (AB).sub.n, fabricated by the Langmuir-Blodgett technique, in the Y-type configuration, comprising n number of bilayers AB. Each bilayer AB consists of a first layer of a sidechain chromophoric polymer A having the electron accepting end of the sidechain chromophore attached to the polymer backbone, and a second layer of another sidechain chromophoric polymer B having the electron donating end of the chromophore attached to the polymer backbone.\nU.S. Pat. No. 5,247,055, issued Sep. 21, 1993 to Stenger-Smith et al, and incorporated herein by reference, describes new syndioregic chromophoric polymers, their self assembly on a surface, and electric field poling to achieve polar films exhibiting second-order nonlinear optical properties. U.S. Pat. No. 5,247,055 teaches the use of a single polymer for all layers. This patent also teaches that when monolayers are built up, one upon the next, by the X or Z types of LB deposition processes, one monolayer may be chemically cross-linked to the next monolayer. For example, carboxyl groups in one monolayer may react with epoxy groups in the next monolayer. Although not discussed in the patent, such chemical cross-linking of adjacent monolayers may override the thermal instability of identical polymer layers due to adjacent hydrophobic/hydrophilic domains and keep their dipole moments pointing in the same direction.\nFor films that are processed by LB or other techniques that align the chromophores by hydrophobic-hydrophilic interactions, the dipole moment of the chromophore has little influence on the actual aligning process (admittedly, dipole-dipole repulsion may have a secondary effect on alignment if the polymer is improperly designed). The largest component of the .beta. tensor generally lies along the long axis of the chromophore, which is also the charge transfer axis of the chromophore. The hydrophilic and hydrophobic groups that force alignment of the asymmetrical chromophore at an interface are positioned at the extreme ends of the long axis of the chromophore. The average direction of alignment of the long axes of the chromophores can be assumed to lie nearly perpendicular to the film surface if hydrophilic-hydrophobic forces are the only forces important to the alignment of the chromophores. However, in LB processing of polymers, there is usually an additional ordering force due to the shearing of the polymer chains as they are pulled from the air-water interface--the polymer backbones tend to align in the dipping direction. In this case, there may also be a preferred azimuthal angle of alignment of the chromophores. The average alignment of the chromophores can be estimated from 3-dimensional birefringence measurements (e.g., from polarized UV-VIS absorption due to the charge transfer oscillation in the chromophore, or from polarized second-harmonic generation measurements).\nThe dipoles of the chromophores do become aligned during LB processing (i.e., they \"go along for the ride\"). Therefore, one can still refer to these films as having \"polar order.\" For the sake of brevity and convenience, in this patent application, the alignment of chromophores in the films will be referred to as imparting \"polar order\" to the films, terminology often found in the literature. As we are interested in the second-order nonlinear optical properties of the films, it will be understood that referring to their \"polar order\" also implies (because of the design of the accordion polymer) that the .chi..sup.(2) tensor of the film is relatively large."}
{"text": "1. Field of the Invention\nThe invention relates to the synthetic methods and manufacture procedures for (methyl-d3)amine and salts thereof.\n2. Description of Related Art\n(Methyl-d3)amine and the hydrochloride thereof, the important intermediates for chemical synthesis, can be used to make medicinal compounds.\nFor example, the ω-diphenylurea derivatives are known as the compounds with c-RAF kinase inhibition activity. Initially, ω-diphenylurea compounds, such as Sorafenib, were firstly found as the inhibitors of c-RAF kinase. The other studies had shown that they could also inhibit the MEK and ERK signal transduction pathways and activities of tyrosine kinases including vascular endothelial growth factor receptor-2 (VEGFR-2), vascular endothelial growth factor receptor-3 (VEGFR-3), and platelet-derived growth factor receptor-β (PDGFR-β) (Curr Pharm Des 2002, 8, 2255-2257). Therefore, it is called multi-kinase inhibitor which possesses dual anti-tumor effects.\nSorafenib (trade name Nexavar), a novel oral multi-kinase inhibitor, was developed by Bayer and Onyx. In December 2005, based on its outstanding performance in phase III clinical trials for treating advanced renal cell carcinoma, Sorafenib was approved by FDA for treating advanced renal cell carcinoma. It was marketed in China in November 2006. However, Sorafenib has various side-effects, such as hypertension, weight loss, rash and so on.\n(Methyl-d3)amine is used during the preparation of sorafenib derivatives. However, the synthetic steps or the current preparation processes are relatively complex, or the cost is high. Therefore, development of some simple, highly efficient, and/or low cost methods for preparing (methyl-d3)amine and salts thereof is needed."}
{"text": "A storage system is known in which a plurality of flash memory packages (hereinafter referred to as “FM packages”) form a parity group in order to enhance redundancy of data. The FM package is a storage device equipped with a plurality of flash memory chips (hereinafter referred to as “FM chips”). PTL 1 describes a storage system that allows the FM package to execute a parity generation process when data related to a write request received from a host is stored in the FM package."}
{"text": "As a disclosure of the NC automatic lathe of this type, there is Patent Document 1, for example. For information, Patent Document 1 belongs to the applicant of this patent application.\nPatent Document 1: JP-A-10-015702\nThe NC automatic lathe having such a configuration as shown in FIG. 7 is disclosed in the above mentioned Patent Document 1. In the first place, there is arranged a headstock 103 provided with a main spindle 101. This headstock 103 is so configured as to move in a direction of Z1 axis which is parallel to an axial direction of the main spindle. A guide bush 105 is arranged in front of the main spindle 101, and holds a tip end portion of a work 107 which is holded by the main spindle 101.\nA back attachment 109 is arranged so as to be opposed to the main spindle 103. This back attachment 109 is provided with a sub spindle 111, and is so configured as to move in a direction of Z3 axis which is parallel to the axial direction of the main spindle 101. A gang tool post 113 is arranged at a side of the guide bush 105, and is so configured as to move in both directions of X1 axis and Y1 axis which are perpendicular to the axial direction of the main spindle 101 and also perpendicular to each other.\nA turret tool post 115 is arranged at another side of the guide bush 105, opposite to the gang tool post 113. This turret tool post 115 is provided with a turret 117, and is so configured as to move in a direction of Z2 axis which is parallel to the axial direction of the main spindle 101, and in both directions of X2 axis and Y2 axis which are perpendicular to the direction of the Z2 axis and also perpendicular to each other.\nThe turret 117 is provided with a plurality of tool holder attaching parts 119, and tool holders which hold tools (not shown) will be detachably attached to the tool holder attaching parts 119. Moreover, a plurality of tools 121 are attached to the gang tool post 113.\nAccording to the above described configuration, front machining is conducted on the work 107 which is holded by the main spindle 101, by means of the tools 121 of the gang tool post 113 and the tools (not shown) attached to the turret tool post 115. After the front machining has been finished, cutting-off work is conducted, and then, back machining is conducted on the work 107 which is holded by the sub spindle 111, by means of the turret tool post 115.\nIn case of the NC automatic lathe having the above described configuration, the back machining on the work 107 has been conducted only by the turret tool post 115 which has been also concerned with the front machining. Accordingly, there has been such a problem that the turret tool post 115 must cover a larger number of machining processes as compared with the gang tool post 113 which is concerned with only the front machining, and it would be impossible to properly distribute the machining processes between the front face and back face thereby to reduce a machining time. Moreover, there has been a problem that because the turret tool post 115 need to be provided with the tools for both the front machining and the back machining, and there is a shortage of tools in number, for conducting complicated and diversified machining. Further, for simultaneously conducting both primary machining and secondary machining, the turret tool post 115 can be concerned with the front machining only with the limited tools, and so, the front machining is conducted only by the gang tool post. In this case, complicated machining cannot be applied even to the front machining, and it would be impossible to reduce the machining time, by simultaneously machining, using the two tool rests.\nUnder the circumstances, an NC automatic lathe having such a configuration as shown in FIG. 8 has been proposed. This is the configuration of the NC automatic lathe of “TNL-26” model, manufactured by Traube GmbH, in Germany.\nIn the first place, there is arranged a headstock 203 provided with a main spindle 201. This headstock 203 is so configured as to move in a direction of Z axis which is parallel to an axial direction of the main spindle 201. Aback attachment 205 is arranged so as to be opposed to the headstock 203. This back attachment 205 is provided with a sub spindle 207, and is so configured as to move in the direction of the Z axis.\nA front machining tool unit 209 is arranged at a side of the back attachment 205. This front machining tool unit 209 is provided with a plurality of front machining tools 211, and is so configured as to move in the direction of the Z axis. Moreover, both the back attachment 205 and the front machining tool unit 209 are so configured as to move also in the direction of X axis.\nA turret tool post 213 and a turret tool post 215 are arranged at both sides of the main spindle 201. The turret tool post 213 is provided with a turret 217, and is so configured as to move in directions of the X axis, Y axis and Z axis. A plurality of tool holders 219 are detachably attached to the turret 217, and tools 221 are attached to the tool holders.\nThe turret tool post 215 is provided with a turret 223, and is so configured as to move in directions of the X axis and Y axis. A plurality of tool holders 225 are detachably attached to the turret 223, and tools 227 are attached to the tool holders.\nMoreover, a back machining tool unit 229 is arranged at a side of the main spindle 201. A plurality of back machining tools 231 are attached to this back machining tool unit 229. The back machining tool unit 229 is so configured as to move in the directions of the X axis and Y axis.\nAccording to the NC automatic lathe having the above described configuration, because there are provided the turret tool post 213 and the turret tool post 215, and additionally, the front machining tool unit 209 is provided, the three tool rests will be simultaneously concerned with the front machining, even though both primary and secondary machining are simultaneously conducted. Accordingly, it is possible to rapidly perform complicated machining on the front face. Further, the back machining is conducted by the back machining tools 231 attached to the back machining tool unit 229 and the back machining tools attached to the turrets 217, 223. As compared with the NC automatic lathes as shown in FIG. 7, the back machining is conducted by the back machining tools 231 attached to the back machining tool unit 229 independently from the front machining, whereby the machining time can be reduced."}
{"text": "In the past, much effort has been devoted to development of fall prevention devices known as rope grips. The desirable qualities of such a device fall into several categories including safety, reliability, versatility and convenience of use. Dodge Machine Company, Inc. commonly offers a rope grip with a number of desirable qualities, for example the rope grip is manufactured from large pieces of material which distributes the impact of a fall to relatively massive pieces of metal providing strength, reliability and safety.\nAnother commercially available rope grip is the Barrow-Hepburn-Everest Rope Grip, patented in the U.K. with Pat. No. 1,077,068. This rope grip operates by means of three small steel balls through which a rope is passed. When the rope grip begins to fall rapidly down a rope, the friction of the rope against the balls draws them into a conical shaped housing which jams the rope between the three balls. While this device offers advantages over many others in fail safe operation, it too has practical drawbacks. For example, this device can only be used with a very narrow tolerance in rope diameter. Moreover, the rope must be threaded through the grip, as it can not be opened for placing the grip on the rope.\nIn U.S. application Ser. No. 501,623 filed on June 6, 1983, a removable double action rope grip is disclosed. This rope grip takes care of many of the practical disadvantages of the prior art, however many problems still remain. In particular, it is often desirable to anchor vertical safety lines on the surface of a nearby structure. According to existing specifications issued by OSHA and other such safety organizations, a safety line should be held in place at least every thirty feet so that it does not blow and flap around in the wind. Whatever means are used to anchor the safety line will, of course, present an obstacle to a rope grip. Conventional safety devices for use on inclined lines cannot pass such an anchor, and therefor the rope grip must be removed from the line and reattached at the desired position.\nThe Latchway.RTM. Cliprider System manufactured and marketed by Barrow-Hepburn Equipment, Ltd. of London, England discloses a safety system that allows workers to move past clips holding horizontal line in place. In addition, this system has special fittings which allow a safety device to be put on or removed from a cable anywhere in the system.\nU.S. Pat. No. 4,265,179 to Tupper describes in detail various aspects of the Cliprider system. By way of summary, the Latchway System is merely a means for protecting a worker by use of a horizontal safety line, and does not provide for arresting movement along the line. That is to say, the worker is protected by a horizontal line while he is working on more or less the same level. This, of course, would not be suitable for use by workers who are climbing long ladders such as found on water storage tanks, television transmitting towers, smokestacks or the like, where a vertical safety line is needed. While the Latchway System does provide a way of bypassing anchors which secure a line, it simply does not act as a safety device for use on an inclined line."}
{"text": "The present invention relates to a multicircuit heat exchanger, i.e. a heat exchanger capable of circulating more than two media. Through delimiting zones in the heat exchanger package, a further medium can be placed in heat exchange with one of the two media normally being circulated. Several such delimited zones may be arranged in the same heat exchanger package, and the zones may even overlap each other for mutual heat exchange.\nIn today\"\"s conventional heat exchangers, two media are circulated in heat exchange with each other. There is, however, a requirement for circulating further media in the exchanger, e.g. for additional heating.\nThe present invention eliminates the above limitation of the prior art by providing delimited zones or islands in the heat exchanger. Inside these zones, a further medium may be circulated and be in heat exchange with one of the other media.\nThe present invention thus provides a multicircuit heat exchanger comprising plates with a pattern of grooves and inlet and outlet connections. The plates are placed in a package and brazed together to form separate ducts for two media between alternate pairs of plates.\nAccording to the invention, at least one delimited zone is defined between alternate pairs of plates, to block off a first medium whereas a second medium is allowed to flow through. The delimited zone has an inlet and an outlet for circulation of a further medium under heat exchange with said second medium.\nThe invention is defined by the accompanying claim 1, whereas advantageous embodiments are stated in the dependent claims."}
{"text": "A plunger lift is an apparatus that is used to increase the productivity of oil and gas wells. In the early stages of a well's life, liquid loading is usually not a problem.\nWhen rates are high, the well liquids are carried out of the tubing by the high velocity gas. As the well declines, a critical velocity is reached below which the heavier liquids do not make it to the surface and start to fall back to the bottom exerting back pressure on the formation, thus loading up the well. A plunger system is a method of unloading gas in high ratio oil wells without interrupting production. In operation, the plunger travels to the bottom of the well where the loading fluid is picked up by the plunger and is brought to the surface removing all liquids in the tubing. The plunger also keeps the tubing free of paraffin, salt or scale build-up. A plunger lift system works by cycling a well open and closed. During the open time a plunger interfaces between a liquid slug and gas. The gas below the plunger will push the plunger and liquid to the surface. This removal of the liquid from the tubing bore allows an additional volume of gas to flow from a producing well. A plunger lift requires sufficient gas presence within the well to be functional in driving the system. Oil wells making no gas are thus not plunger lift candidates.\nAs the flow rate and pressures decline in a well, lifting efficiency declines geometrically. Before long the well begins to “load up”. This is a condition whereby the gas being produced by the formation can no longer carry the liquid being produced to the surface. There are two reasons this occurs. First, as liquid comes in contact with the wall of the production string of tubing, friction occurs. The velocity of the liquid is slowed, and some of the liquid adheres to the tubing wall, creating a film of liquid on the tubing wall. This liquid does not reach the surface. Secondly, as the flow velocity continues to slow the gas phase can no longer support liquid in either slug form or droplet form. This liquid along with the liquid film on the sides of the tubing begin to fall back to the bottom of the well. In a very aggravated situation, there will be liquid in the bottom of the well with only a small amount of gas being produced at the surface. The produced gas must bubble through the liquid at the bottom of the well and then flow to the surface. Because of the low velocity very little liquid, if any, is carried to the surface by the gas. Thus, as explained previously, a plunger lift will act to remove the accumulated liquid.\nA typical installation plunger lift system 100 can be seen in FIG. 1 (prior art). Lubricator assembly 10 is one of the most important components of plunger system 100. Lubricator assembly 10 includes cap 1, integral top bumper spring 2, striking pad 3, and extracting rod 4. Extracting rod 4 may or may not be employed depending on the plunger type. Below lubricator 10 is plunger auto catching device 5 and plunger sensing device 6. Sensing device 6 sends a signal to surface controller 15 upon plunger 200 arrival at the well top. Plunger 200 is shown to represent the plunger of the present invention and will be described below in more detail. Sensing the plunger is used as a programming input to achieve the desired well production, flow times and wellhead operating pressures. Master valve 7 should be sized correctly for tubing 9 and plunger 200. An incorrectly sized master valve will not allow plunger 200 to pass. Master valve 7 should incorporate a full bore opening equal to the tubing 9 size. An oversized valve will allow gas to bypass the plunger causing it to stall in the valve. If the plunger is to be used in a well with relatively high formation pressures, care must be taken to balance tubing 9 size with the casing 8 size. The bottom of a well is typically equipped with a seating nipple/tubing stop 12. Spring standing valve/bottom hole bumper assembly 11 is located near the tubing bottom. The bumper spring is located above the standing valve and can be manufactured as an integral part of the standing valve or as a separate component of the plunger system.\nSurface control equipment usually consists of motor valve(s) 14, sensors 6, pressure recorders 16, etc., and electronic controller 15 which opens and closes the well at the surface. Well flow ‘F’ proceeds downstream when surface controller 15 opens well head flow valves. Controllers operate on time, or pressure, to open or close the surface valves based on operator-determined requirements for production. Modern electronic controllers incorporate features that are user friendly, easy to program, addressing the shortcomings of mechanical controllers and early electronic controllers. Additional features include battery life extension through solar panel recharging, computer memory program retention in the event of battery failure, and built-in lightning protection. For complex operating conditions, controllers can be purchased that have multiple valve capability to fully automate the production process.\nIn these and other wells it is desirable to measure the downhole temperature and pressure versus time, chemical profiles and other data. This information is used to figure oil and gas reserves and production plans. Conventional methods include dropping special sensors called pressure bombs via cable down the tubing. Pressure bombs can be attached to the wireline or left downhole to be retrieved by fishing at a later date. Special trucks with a crew are used which is expensive for the well operator.\nIn FIG. 2 (prior art), a special truck called a wireline (also called slickline) rig 200 is used to drop a downhole equipment data logger (temperature and/or pressure and/or time) 207 down tubing 266 of the well. Nominally the tubing is two inches in diameter, and data logger 207 is about three feet long. Wireline rig 200 has an on-board computer 201 for data recording. Hoistable crane 202 supports electric line 206 which usually requires a lubricator 203 and a blowout protector 204. A spool and hoist assembly 205 controls electric line 206. All this special equipment is costly to lease for the well operator. Furthermore, the use of this equipment requires the complete shutdown of the well during the operation of dropping special data logger 207.\nWhat is needed is an improved data logger sensor that can be dropped down a well and retrieved without a wireline rig. The plunger will house and deliver the data logger to the bottom of the well to take readings. Then the well operator can turn the well on to flow the plunger and data logger to the surface without the use of a wireline rig and crew. This sensor should be easily detachable to the plunger and readily plugged into a computer to retrieve the measured downhole temperature and/or pressure. The present invention fulfills these needs for the well operator/producer."}
{"text": "A physician who performs an operation on the body of a patient, generally employs a real-time imaging system, in order to view the location and orientation of the medical intervention device (e.g., catheter, needle), within the body of the patient during the operation. Such real-time imaging systems are known in the art. These systems generally enable a display to display a representation of the medical intervention device superimposed on an image of the body of the patient.\nU.S. Pat. No. 6,351,513 issued to Bani-Hashemi et al., and entitled “Fluoroscopy Based 3-D Neural Navigation Based on Co-Registration of Other Modalities with 3-D Angiography Reconstruction Data”, is directed to a method for displaying a real-time 3-D reconstruction of a catheter within a 3-D angiography reconstruction of a vessel. The method includes the procedures of acquiring a 3-D angiography image of the arterial tree by a computed tomography device and registering the 3-D angiography image with a 2-D fluoroscopic image of a vessel, according to the structural similarities (i.e., anatomical landmarks).\nThe method further includes the procedures of determining the projecting lines of the catheter by using an X-ray apparatus, determining the location of the catheter, by intersecting the 3-D angiography image with the projecting lines and displaying a 3-D visualization of the 3-D reconstruction of the catheter within the 3-D angiography reconstruction of the vessel. The 3-D visualization of the catheter is updated as the catheter is moved.\nU.S. Pat. No. 6,314,310 issued to Ben-Haim et al., and entitled “X-Ray Guided Surgical Location System with Extended Mapping Volume”, is directed to a system for inserting a needle into a selected location of the vertebrae of a patient. The system includes a reference element, a plurality of magnetic field generator coils, a driver circuitry, a computer, a user interface control, a display, a fluoroscope and a computer tomography (CT) device. The reference element is in form of a plastic disc transparent to visible light and X-rays, which includes three equally spaced metal fiducial marks at the periphery thereof, a first position and orientation sensing device at the center thereof and another fiducial mark adjacent the first position and orientation sensing device. The needle includes a second position and orientation sensing device.\nThe magnetic field generator coils are placed on or adjacent to a bed on which the patient lies. The fluoroscope irradiates the patient from one side of the body of the patient. The computer controls multiple aspects of the system. The first position and orientation device and the second position and orientation device sends signals to the computer, respective of the time-varying magnetic fields generated by the magnetic field generator coils. The computer analyzes the signals to determine the six-dimensional position and orientation coordinates of the first position and orientation device and the second position and orientation device, relative to a common frame of reference defined by the magnetic field generator coils. The computer enables the display to display an image of the vertebrae, a representation of the first position and orientation device and the second position and orientation device and a representation of the needle and the fiducial marks. The location and the angular orientation of the reference element are determined by determining the two-dimensional coordinates of the representation of the fiducial marks. A scaling factor is determined for the images displayed on the display, by comparing the determined coordinates with the known positions of the fiducial marks.\nWhile acquiring CT images of the body of the patient, the reference element is fixed to the body and remains fixed to the body in this position during the surgery. The CT images are registered with the X-ray images, by comparing the image-derived coordinates of the fiducial marks of the reference element, which appear in the CT images, with the image-derived coordinates of the fiducial marks in the X-ray images. The fiducial marks of the reference element and the fiducial marks in the X-ray images are visible marks. The three-dimensional CT images are rotated or scaled, in order to align the CT images with the X-ray images and the CT images are projected onto the plane of the X-ray images and superimposed on the X-ray images or displayed alongside the X-ray images.\nU.S. Pat. No. 6,421,551 issued to Kuth et al., and entitled “Method for Registering Images of a Subject with a Magnetic Resonance System and Magnetic Resonance System for the Implementation of the Method”, is directed to a system for readjusting the tomogram plane of an image of the body of a patient. The system includes a control console, a magnetic resonance system, a stereoscopic camera and a marking element. The control console includes a control unit, an image data generator and processor, a coordinate transformation unit, a readjustment unit and a tomogram selecting unit. The magnetic resonance system includes two pole shoes which are located opposite one another.\nThe control console is connected to the magnetic resonance system and to the stereoscopic camera. The marking element is composed of three reflective balls and is arranged at the patient in the region of the knee joint, in a first coordinate system. The stereoscopic camera acquires an image of the reflective balls and sends the respective position data to the control console. The coordinate transformation unit transforms the position data from the first coordinate system to a second coordinate system of the magnetic resonance system. When the relative movement of the patient is known, the readjustment unit readjusts the previously defined tomogram plane, such that it again lies relative to the marking element with respect to the knee joint, as it did in the preceding joint position.\nOne way to destroy tumors in a patient, and to prevent metastasis, is by subjecting the target tissue to radiation therapy. One type of radiation therapy is known as linear acceleration, whereby a beam of x-rays or electrons is directed at the target tissue from different directions. Each time the linear accelerator directs a beam towards the target tissue it also irradiates healthy tissue which surrounds the target tissue, along the path of the irradiation beam. Accordingly, such surrounding tissue is irradiated significantly less than the target tissue.\nThe linear accelerator is programmed to irradiate a specific volume which is generally similar to the shape of the target tissue. Accordingly, the portion of the body including the target tissue, has to be placed such that the target tissue is located within that specific volume. A conventional linear acceleration treatment includes a plurality of recurring procedures, usually over a period of several days or weeks. Each time, the portion of the body including the target tissue, has to be placed exactly as it was placed in the first treatment.\nFor this purpose, during the first radiation session, after locating the portion of the body which contains the target tissue at a location appropriate for irradiation, a plurality of non-hazardous laser beams, for example four beams, are directed from fixed locations, toward that portion of the body. These four points are marked by a permanent marker, such as a waterproof marker, on the skin of the patient. At every subsequent session, that portion of the body is re-positioned to the position and orientation determined at the first session, by directing the same four laser beams toward the same portion of the body and repositioning that portion, until the four permanent marks line up with the four laser beams."}
{"text": "Wireless communication systems are widely deployed to provide various telecommunication services, such as telephony, video, data, messaging, and broadcasts. Typical wireless communication systems may employ multiple-access technologies capable of supporting communication with multiple users by sharing available system resources (e.g., bandwidth, transmit power, etc.). Examples of such multiple-access technologies include code division multiple access (CDMA) systems, time division multiple access (TDMA) systems, frequency division multiple access (FDMA) systems, orthogonal frequency division multiple access (OFDMA) systems, single-carrier frequency divisional multiple access (SC-FDMA) systems, and time division synchronous code division multiple access (TD-SCDMA) systems.\nThese multiple access technologies have been adopted in various telecommunication standards to provide a common protocol that enables different wireless devices to communicate on a municipal, a national, a regional, and even a global level. An example of a telecommunication standard is Long Term Evolution (LTE). LTE is a set of enhancements to the Universal Mobile Telecommunications System (UMTS) mobile standard promulgated by Third Generation Partnership Project (3GPP). LTE is designed to better support mobile broadband Internet access by improving spectral efficiency, lowering costs, improving services, using new spectrum, and integrating with other open standards using OFDMA on the downlink (DL), SC-FDMA on the uplink (UL), and multiple-input multiple-output (MIMO) antenna technology."}
{"text": "1. Field of the Invention\nThe present invention relates to a pilot operated electromagnetic valve.\n2. Description of the Related Art\nAn automotive air conditioner is generally configured such that it includes a compressor, a condenser, an evaporator, and so forth arranged in a refrigerant circulation passage. Various types of control valves are provided for the purpose of, for example, switching the refrigerant circulation passages according to the operation state in such a refrigeration cycle and regulating the flow rate of refrigerant. A pilot operated electromagnetic valve capable of controlling the opening and closing of a large valve section using a relatively small electric power may be used as such the control valve (see Reference (1) in the following Related Art List, for instance).\nSuch an electromagnetic valve drives a small pilot valve element by a solenoid so as to open and close a pilot valve and then drives a large main valve element by a pressure difference regulated thereby so as to open and close a main valve. A piston is formed integrally with the main valve element, and a back pressure chamber is formed by this piston as a separated space inside a body. A leak passage, through which to introduce the refrigerant into the back pressure chamber and a pilot passage, through which the refrigerant is led out from the back pressure chamber, are formed in the main valve element. And the opening and closing of the pilot valve opens and blocks the pilot passage, respectively. With this structure, the opening and closing of the main valve is controlled by varying the pressure of the back pressure chamber. The pressure of the back pressure chamber is regulated through a balance between the flow rate of refrigerant introduced into the back pressure chamber and the flow rate of refrigerant led out from the back pressure chamber."}
{"text": "The present invention relates generally to electrochemical machining, and, more specifically, to machining of long tubes.\nPetroleum undergoes pyrolysis in cracking furnaces for producing useful products such as ethylene. The cracking furnace uses long tubes through which the petroleum is channeled and heated for the pyrolysis thereof.\nFluted process tubes are desired for increasing the available surface area and the corresponding rate of pyrolysis. Such fluted tubes may be formed by extruding suitably malleable metal, yet such metal is subject to undesirable creep at the high temperature required in the pyrolysis cracking process.\nAccordingly, a nickel and chromium based superalloy having substantial high temperature creep strength is preferred for the process tubes. However, such high strength material cannot be extruded, and the flutes must be otherwise formed.\nFor example, electrochemical machining (ECM) may use a cathode tool in the desired shape of the internal flutes of the tube for electrochemically forming the flutes as an electrolyte flows between the tool and the tube when powered under high electrical current defining respective cathode and anode. ECM chemically erodes the initially circular unfinished surface defining the tube bore to form internal flutes being complementary with the external flutes of the cathode tool.\nHowever, since the individual tubes are long, and the cathode tool must be driven completely through the tube, the electrolyte contained in the tube can effect stray ECM in the finished bore as the tool continues its travel through the tube. And, the electrolyte becomes heated during operation and affects the ECM process.\nAccordingly, it is desired to provide an improved ECM apparatus and process for locally isolating the cathode tool during operation, as well as providing internal cooling of the tube in a preferred embodiment.\nAn ECM apparatus includes a fluted cathode tool for being driven through a tubular workpiece for electrochemically forming internal flutes therein. The cathode tool is sealed at one end for isolating electrolyte over the flutes thereof to seal against flow of the electrolyte past the tool and into the finished fluted bore."}
{"text": "An infusion pump assembly may be used to infuse a fluid (e.g., a medication or nutrient) into a user. The fluid may be infused intravenously (i.e., into a vein), subcutaneously (i.e., into the skin), arterially (i.e., into an artery), and epidurally (i.e., into the epidural space).\nInfusion pump assemblies may administer fluids in ways that would be impractically expensive/unreliable if performed manually by nursing staff. For example, an infusion pump assembly may repeatedly administer small quantities of an infusible fluid (e.g., 0.1 mL per hour), while allowing the user to request one-time larger “bolus” doses."}
{"text": "In wireless communication systems, effort has been made to increase the use of spectrum to allow for a greater number of users of a given frequency band. One example of a technique to increase spectrum efficiency is a frequency division multiple access (FDMA) technique. In a conventional FDMA system, a given frequency band is divided into a number of channels, wherein each channel is occupied by one user. An FDMA system can also be a time division duplex (TDD) system wherein a given RF channel is used for transmitting frames in both the forward and reverse directions of communication which are separated in time.\nOther techniques comprise digital multiple access communication systems. One such conventional digital multiple access technique for increasing efficient use of spectrum is a time division multiple access (TDMA) technique. In a TDMA system, each channel is divided into a plurality of slots for the transmission of signals. Each time slot may be allocated to a different call. A TDMA system can also employ TDD techniques. Accordingly, a number of calls can be simultaneously transmitted on a single channel or frequency.\nFinally, increased spectrum efficiency can be achieved by spread spectrum techniques, in the form of either a slow frequency hopper system or a direct-sequence CDMA system. In a slow frequency hopper system, the carrier frequency of the signal is changed at a predetermined rate over a wide range of possible frequencies in a pseudo-random sequence known in advance by the receiver. Generally, spread spectrum techniques reduce the effects both intentional or unintentional interference. Direct sequence CDMA systems allow multiple users to share the same spectrum wherein each user is assigned a unique pseudonoise code sequence. The signal is spread by the wide bandwidth pseudo-noise sequence known in advance by the receiver.\nIn a conventional multiple access communication system, a base station periodically transmits a beacon or synchronization signal to maintain synchronization with an associated handset. For example, in a time division duplex multiple access system, a slot in a time frame must be occupied even if the base station is not in a call. Depending upon the communication system, available resources for providing communication between devices are used to maintain synchronization between devices which are not in communication. Accordingly, there is a need for a method and apparatus for increasing spectrum efficiency to allow a greater number of users to occupy a given frequency band."}
{"text": "The goal of vegetable breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits may include any trait deemed beneficial by a grower and/or consumer, including greater yield, resistance to insects or disease, tolerance to environmental stress, and nutritional value.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant or plant variety. A plant cross-pollinates if pollen comes to it from a flower of a different plant variety.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants of different genotypes produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform varieties requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods that have been used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more plants or various other broad-based sources into breeding pools from which new lines and hybrids derived therefrom are developed by selfing and selection of desired phenotypes. The new lines and hybrids are evaluated to determine which of those have commercial potential."}
{"text": "This invention relates to a respiration monitor, in which the respiration activity is monitored by measuring the variations of the thorax impedance or other suitable parameters of a patient caused by his respiration activity.\nA sufficient gas exchange in the lung is assured only if the lung is sufficiently ventilated. This ventilation is dependent on the respiration frequency on the one hand and on the respiration \"deepness\" (breath volume) on the other hand. Respiration monitors are used for monitoring the respiration of patients exposed to danger, e.g., patients in intensive care stations or prematurely born children. Generally, such instruments register the respiration activity and produce an alarm signal in the case of cessation of respiration (apnea), if the apnea duration exceeds a predetermined alarm limit.\nIn most cases the respiration monitoring is based on an impedance measurement of the patient's thorax. However, other measuring methods are also known, e.g., measuring the temperature variations of the respiration air or pressure variations between body and support.\nEffective respiration is possible with relatively low frequency and deep breathing as well as with relatively high frequency and flatter breathing. Some prior art respiration monitors comprise an automatic sensitivity control responsive to the amplitude of the respiration signals. This control is performed with a given delay corresponding to a given time constant. In order to prevent misinterpretation of noise and artifacts as respiration signals, a minimum amplitude threshold is provided, below which respiration signals are no longer detected. If this threshold is not reached for a duration exceeding the alarm limit an alarm signal is produced.\nSince the minimum amplitude threshold must be so low that all effective breaths are detected, it may happen that spurious signals exceed this threshold and are misinterpreted as respiration signals although a possible apnea is present. Such spurious signals may have their origin in the heart activity, for example. A respiration monitor capable of suppressing such spurious signals is described in U.S. Pat. No. 3,976,052, issued to Gerhard Junginger and Helmut Zeeb and assigned to the assignee of this application. The respiration monitor described in this patent operates with two amplitude thresholds, one of which is a fixed minimum threshold and the other one of which is a variable threshold dependent on the respiration amplitude. The respiration monitor is connected to separately receive signals corresponding to the heart beat rate and compare the periods of these signals to the periods of the signals measured at the thorax of the patient. If the periods of both signals are substantially equal, the signals measured at the thorax are actually due to heart activity rather than to respiration. In this case, the variable threshold is raised so that the monitor no longer responds to heart beat related signals but indicates that no respiration signals are measured.\nAlthough capable of good suppression of spurious signals, prior art respiration monitors do not take into account the fact that certain further conditions must be met in order to assure a gas exchange that is sufficient for the patient. For example, if the respiration frequency is low, a predetermined minimum deepness of breath must be achieved."}
{"text": "A boat can be more easily maneuvered by the use of a bow thruster which consists of a water pump and a valve arrangement for diverting the water into either of two pipelines that open at the opposite sides of the boat near the bow. It would be desirable if the thrust obtainable from a water pump placed on board a vessel could be used to help maneuver the boat, or if required, to propel the boat forwardly, as when the propeller has been broken or where the propeller could cause injury to persons or articles in the water."}
{"text": "1. Field of the Invention\nThe present invention relates to a network game in which a plurality of players participate to progress with the game, and, more particularly, to execution of an event in which a plurality of players participate.\n2. Description of the Related Art\nThe recent advancement on the network technology has made network games (online games) vivid in each of which a plurality of players participate to progress with the game. A system that executes a network game includes a server apparatus which is managed by one who runs the network game, and client devices of individual players which are connectable to the server apparatus over a network such as the Internet.\nAn RPG (Role Playing Game) among the network games progresses as multiple players participate in the game at the same time so that their player characters cooperate with one another to battle with an opponent character, or the players make their player characters battle with one another. Recently, in particular, an MMORPG (Massively Multiplayer Online RPG) which permits participation of more players is becoming popular. A network game such as an MMORPG requires that the movements of the player characters of all the participating players should be transmitted to the client devices manipulated by the players without contradiction.\nJapanese Patent Application Laid-Open Publication No. 2008-99906 discloses a technique of allowing a server apparatus to centralize management on the timings for the actions and processing of the player characters, and periodically synchronize with the timing of progressing the game. The technique disclosed in Japanese Patent Application Laid-Open Publication No. 2008-99906 prevents the individual client devices from processing time-sequential information with some delays which would bring about inconsistency on the progress of the game. This scheme can be adapted to processing of an event in which a plurality of players participate in a network game.\nWhen an event occurs in a network game, a movie may be played back on the individual client devices as a previous step to the event so that the individual players watch the movie at the same time. In a network game such as an MMORPG in which multiple players participate, however, there are a variety of players including skilled or experienced players and beginners. In a case where experienced players and beginners participate in the same event, is likely that the experienced players have already seen such a movie so that they may feel tedious about being tied on watching the movie for a long time whereas the beginners are apt to see it for the first time.\nFrom the viewpoint of just playing back a movie, the each client device may be permitted to freely skip the movie as each player is likely to skip some scenes in progressing with a stand-alone game, so that the player is relieved of a long tie-up. When an event such as a battle immediately follows the end of a movie, however, the timing for a process related to the event in the progress of the game may vary among the individual client devices. From the viewpoint of playing an event such as a battle with a plurality of players in cooperation, it is not desirable that while some players are watching a movie, the other players freely start the battle."}
{"text": "The present invention relates generally to methods and apparatuses for computer-based education.\nEducational software that tests a user's knowledge, coupled with the immediate scoring of answers, is well known in the art. For home use, a wide range of software is directed to various age groups including school-grade children. Thus, many educational programs are designed to be quite entertaining, with built-in tests being played as games having an underlying educational purpose. Thus, terms such as “student”, “player”, or “test-taker,” shall be understood to mean any participant in, or user of, educational software that tests the user's knowledge in a formal or informal manner.\nSoftware which blends both fun and learning is often referred to as “edutainment” software. Some edutainment programs include “Mathblaster”, “Where in the World is Carmen Sandiego,” and “Word Munchers.” These edutainment programs (as well as certain other educational programs generally) present players with a series of increasingly difficult tests or puzzles, wherein players must correctly solve the present round before they are allowed to continue to the next round of play.\nIn the above-mentioned edutainment software, each player exists in isolation from all other players. This arrangement is reasonably well suited to informal, isolated learning, but lacks the comprehensive data collection/comparison features needed for formal test administration.\nAnother drawback of the existing home education software is that standardized entertainment elements often become repetitive and boring after seeing them for the third or fourth time. Users often stop using the software once the novelty has worn off."}
{"text": "1. Field of the Invention\nThis invention relates to scroll fluid machine, in which sucked fluid is compressed with stationary and revolving scrolls and discharged to the outside.\n2. Description of the Related Art\nA scroll fluid machine compresses fluid sucked from its peripheral part in a sealed space formed by its stationary and revolving scrolls progressively as the fluid is fed toward its central part, and discharges the compressed fluid from the central part. As the fluid is compressed, the temperature in the sealed space formed by the wraps is elevated. This poses a problem that bearings, seal members, etc. provided in drive parts are soon deteriorated. Heretofore, the scrolls are cooled to hold the temperature within a predetermined temperature.\nWell-known cooling systems cool either a non-driven part, i.e., the stationary scroll, or a driven part, i.e., the revolving scroll.\nFIG. 16 shows a technique concerning a non-driven part cooling system. As shown, a revolving scroll 116 which is mounted on a frame 109 provided in a sealed housing 105, comprises a disc-like body 114 having a shaft 113 depending therefrom. The frame 109 has a central hole, in which a drive shaft 104 coupled to a drive (not shown) is fitted for rotation, and the shaft 113 is eccentrically coupled to the drive shaft 104. The revolving scroll 116 has a wrap 115 engaging with a wrap 111 of a stationary scroll 112.\nThe stationary scroll 112 has a peripheral wall having a suction hole 118. When the revolving scroll 116 is revolved relative to the stationary scroll 112 with the rotation of the drive shaft 104, a sealed space formed by the wraps 111 and 115 is progressively reduced in volume, thus compressing gas entering the sealed space. The compressed gas is discharged from a discharge hole 121 formed in a central part of the stationary scroll 112 through a discharge pipe 120 to the outside.\nA plurality of radially spaced-apart heat pipes 122 are provided in the body 110 of the stationary scroll 112 to remove heat generated in a compression stroke as described above.\nFIG. 17 shows a well-known cooling system for cooling driven part, i.e., the revolving scroll.\nA housing 211 as shown comprises a rear and a front housing part 212 and 213, and a drive shaft 214 is supported for rotation by bearings 215 in a bearing portion of the rear housing part 212. The drive shaft 214 has an extension projecting outward from the bearing portion and coupled to a motor (not shown). The drive shaft 214 also has an eccentric portion 214b, which has an eccentric axis 02--02 with respect to the axis 01--01 of the drive shaft 214 by a distance .delta..\nA revolving scroll 216 which is coupled to the eccentric portion 214b of the drive shaft 214, has a disc-like plate 216a having a mirror finished front surface, a spiral wrap 216b formed on the front side of the mirror finished plate 216a, a boss 216c formed as the driving center with an axial line 02--02 on the rear side of the plate 216a and having smaller diameter than the inner peripheral surface edge of above portion 213b, a ring-like ridge 216d formed on the rear side of the above 216a and on the periphery thereof, and a plurality of radial vent holes 216e formed in a diameter direction above the ridge 216d.\nA stationary scroll 221, which is secured to the front housing part 213, has a disc-like plate 211a having a mirror finished rear surface, a spiral wrap 221b formed on the rear side of the plate 211a and a peripheral wall 221c surrounding the wrap 221b.\nThe wraps 216b and 221b of the revolving and stationary scrolls 216 and 221 engage with or wrap each other at a predetermined deviation angle, and they form a plurality of compression chambers or spaces when the revolving scroll 216 is revolved.\nThe drive shaft 214 has a counterweight 225 mounted on its portion extending in the rear housing part 212, and a centrifugal fan 226 is mounted on the counterweight 225 to generate cooling air flow with the rotation of the drive shaft 214.\nIn the prior art non-driven part cooling system shown in FIG. 16, in which the heat pipes 122 are provided in the stationary scroll body, the heat absorbing portions of the heat pipes 122 are more remote from the revolving scroll which is driven than from the stationary scroll. Therefore, the neighborhood of the bearings, seal members and other parts which are driven in contact with the revolving scroll 116 in the driving thereof, is cooled less efficiently compared to the cooling of the stationary scroll. This means that uniform temperature distribution cannot be obtained.\nThe heat radiating portions of the heat pipes 122 are cooled by their heat radiation to the sealed housing inner space 105a, which is filled with gas-sucked through a suction pipe 119.\nIn communication with the space 105a is the suction hole 118, through which gas enters the compression space which is formed by the stationary and revolving scrolls. This means that gas having been elevated in temperature by the heat radiation from the heat pipes 122 again enters the compression space through the suction hole 118, thus reducing the cooling efficiency.\nIn order to prevent the cooling efficiency reduction, it is necessary to provide special cooling means on an external part to which the suction pipe 119 is connected, thus complicating the construction and increasing the size of the apparatus.\nIn the well-known driven part cooling system shown in FIG. 17, with the rotation of the drive shaft 214 external gas is sucked through a suction passage 227 by the centrifugal fan 226 and led through a ring-like space B and a cooling air passage 220 to be discharged through a discharge passage 228.\nSince in this system the gas having cooled down a central part of the revolving scroll 216 is discharged along the rear side of the revolving scroll 216 and through the discharge passage 228, the provision of the discharge passage is necessary. In addition, in order to increase the cooling efficiency, a cooling fan for cooling the rear side of the stationary scroll 221 has to be provided, thus increasing the size of the apparatus."}
{"text": "In the formation of a double lock stitch, the thread carried by the sewing machine needle through the opening in the stitch plate is engaged below the stitch plate by a gripper which can, in some machines, be of the rotatable type. A double-lock-stitch sewing machine with rotating gripper generally comprises, with the gripper, a bobbin housing from which the bobbin feeds a looper thread or underthread for engagement with the thread carried downwardly by the needle to produce the double lock stitch.\nIn general, the bobbin housing is provided with a holding device adapted to prevent the bobbin housing from being frictionally entrained by the rotation of the gripper.\nSuch holding device can comprise a holding finger, on the one hand, and an abutment on the other hand. During each stitch formation in the conventional system, a throughgoing gap is opened between the retaining finger and the abutment to pass the upper thread loop. This gap formation can be effected by the thread itself, in which case, under the effect of thread tension, the bobbin housing is rotated slightly in the sense opposite the gripper rotation sense through a small angle. This results in the gap between the retaining finger and the abutment to clear the upper thread loop.\nWith rapidly operating industrial sewing machines, however, the force which is effective on the holding device is generally so great that the upper thread tension must be increased to an undesirable extent to bring about the aforementioned gap formation.\nIn one prior art double-lock-stitch sewing machine (see German patent application No. 14 39 214), the holding device for the bobbing housing comprises a two-part retaining finger with a rigid and a yieldable section. The yieldable portion of the retaining device is formed as a spring rod which is movable in the direction of withdrawal of the needle-thread loop and can be resiliently deformed by the thread to permit the aforementioned throughgoing gap to form.\nSince the yieldable mass is small, a passage for the needle-thread loop is achieved without interfering with the formation of clean seam or undesirable increase in the thread tension.\nHowever, the retaining device for the bobbin housing of this prior art teaching is suitable only for a gripper rotatable in a vertical plane and makes use of a two-part retaining finger which is fixed to the fabric-carrying plate and cooperates with abutments on the bobbin housing. Such a retaining device has been found to be unsatisfactory for grippers rotatable in a horizontal plane because of the different thread movements and the inability to mount the retaining finger and abutments in the aforedescribed manner."}
{"text": "The present invention relates to a new and distinct cultivar of Nemesia, botanically known as Nemesia caerulea×Nemesia strumosa and hereinafter referred to by the name ‘Kirine-4’.\nThe new Nemesia is 2 product of a planned breeding program conducted by the Inventor in Tochigi, Japan. The objective of the breeding program is to create new Nemesia cultivars with good vigor and attractive flower coloration.\nThe new Nemesia originated from a cross-pollination made by the Inventor in March, 2004 in Tochigi, Japan of the Nemesia caerulea cultivar White Wings, not patented, as the female, or seed, parent with the Nemesia strumosa cultivar Nebula White, not patented, as the male, or pollen, parent. The new Nemesia was discovered and selected by the Inventor as a single flowering plant within the progeny of the stated cross-pollination in a controlled environment in Tochigi, Japan in September, 2004.\nAsexual reproduction of the new Nemesia by terminal cuttings in a controlled environment in Tochigi, Japan since October, 2004, has shown that the unique features of this new Nemesia are stable and reproduced true to type in successive generations."}
{"text": "Adhesively backed index tabs are known in the art. The tabs typically have a coating of removable or permanent pressure sensitive adhesive on the back, such that the tab may be placed on the edge of the page or other substrate. The tabs may be used, for example, to mark a page in a book. A tab may have a portion that extends beyond the edge of the page, to identify the page that is marked. The tab may have a writable portion, such that the user can write or print text and/or graphics.\nIndex tabs of this type are typically supplied on a sheet, from which the user may remove the index tabs as desired. The sheets are typically U.S. letter size or similar and are not always convenient to store due to the relatively large size. Alternatively the tabs may be supplied affixed to sheets of a pad. A sheet of tabs that is removed from a pad is often easily misplaced because it is thin. Alternatively, the tabs may be supplied in a dispenser, with the user pulling the tabs from the dispenser as needed. The dispensers are typically small and easily misplaced or somewhat bulky such that they cannot be conveniently kept where most needed."}
{"text": "A laser welding has been used for manufacturing a composite product which comprises a plurality of resin parts.\nJP-2002-248688A shows that one resin part with a convex portion is welded to the other resin part by a laser beam, while the convex portion is pressed to the other resin part. JP-2003-19752A shows that one resin part is press-inserted to the other resin part receiving a laser beam so that two resin parts are welded together.\nHowever, in JP-2002-248688A, since an end surface of the convex portion and a contacting surface of the other resin part have surface roughness, some clearances are not avoidable therebetween. Due to these clearances, the resin parts are not well welded to generate voids in a welded contacting portion, so that the strength and air-tightness of the welded contacting surface are deteriorated.\nIn JP-2003-19752A, when the resin part has a complicated shape or when the insert amount is small, some clearances are generated between contacting surfaces of two resin parts, whereby the resin parts are not well welded to generate voids in a welded contacting surface, so that the strength and air-tightness of the welded contacting surface are deteriorated. Even if the shape of the resin parts is simplified and the insert amount is made large, some cracks may be arisen in the resin parts."}
{"text": "1. Field of the Invention\nThis invention relates to a color image processing technique for converting an input color image signal into a color image record signal of a color printer, etc., and more particularly to a technique for converting an input color image signal into a color image record signal of N or more colors including a solid color (N is an integer of five or more).\n2. Description of the Related Art\nIn a print technique, solid colors provided by primary color ink of red (R), green (G), and blue (B) and fluorescent ink are added to four colors of Y (yellow), M (magenta), C (cyan), and K (black) to perform color reproduction as a technique for representing vivid colors that cannot be reproduced in four-color process print of Y, M, C, and K. As color samples of solid colors, color samples of Pantone and the like are known and about 1000 solid colors are defined.\nHowever, up to a maximum of eight colors can be printed at the same time on a printer and thus the number of solid colors that can be used is four at the maximum and is a very small number; there is designer's demand for increasing the number of solid colors that can be used. Further, to prepare solid color ink, basic inks of about 20 colors need to be mixed in accordance with the mixing ratio defined in the color sample. When solid color ink is replaced, ink mixing work, printer cleaning work, and the like are required and thus solid color print involves a large number of steps at the print site; this is a problem.\nTo solve such a problem, in print technique, a color reproduction method called HiFi color print is proposed for the purpose of enlarging color gamut by perform color reproduction using five to seven colors in which at least one color of RGB is added to four process colors of YMCK, and making it possible to reproduce solid color colorimetrically by performing color separation processing. As HiFi color print, hexachrome print of Pantone is widely used; it is known that about 90% of solid color can be reproduced by performing color reproduction in six colors of four colors of YMCK plus R-based orange (O) ink and G ink. The color reproduction method of the HiFi color print, which will be hereinafter referred to as HiFi color, is general as a technique of enlarging the color gamut since the application range of HiFi color is also widened to ink jet and electrophotographic color printers.\nFurther, in print and a color printer, a color signal of a display or a digital camera of sRGB color space, etc., may be input as an input signal. Since the color gamut of YMCK four-color process print or a color printer is narrower than the color gamut of the sRGB color space, there is demand for enlarging the color gamut intended for faithful reproduction. Also in such a case, enlarging the color gamut using the HiFi color is effective.\nThus, in the HiFi color, it is possible to enlarge the color gamut of a color printer. At this time, to well reproduce a calorimetric value signal of input solid color and an sRGB color signal, it is necessary to match the tristimulus values of print of output matter with the tristimulus values of the input color signal, namely, realize colorimetric color reproduction in color conversion processing of converting the input color signal into a an image record signal of five to seven colors.\nA color management system (CMS) based on the specifications proposed by International Color Consortium (ICC) widely prevalent as industrial standard at present converts an RGB color signal of an input color signal of machine-dependent color space into a device-independent color signal like L*a*b* or XYZ color space and then converts the device-independent color signal into a YMCK color signal of a color printer of machine-dependent color space. As such conversion processing is performed, the input color signal and the output image record signal match indevice-independent color signal and thus it is possible to ensure colorimetric color reproduction. As the CMS, ColorSync installed in the Mac (registered trademark) operating system of Apple and ICM installed in Windows (registered trademark) of Microsoft are representative.\nIt is also considered in the HiFi color that a CMS in conformity with ICC is adopted for realizing colorimetric color reproduction. To do this, color conversion from a device-independent color signal like L*a*b* or XYZ color space to an image record signal of five to seven colors maybe realized. The following technique is proposed as a related art of color conversion processing for the HiFi color:\nU.S. Pat. No. 4,812,899 discloses an image formation apparatus for performing color reproduction in process inks of seven colors of YMCKRGB. In the image formation apparatus, color conversion from an RGB signal of an input machine such as a scanner to YMCKRGB of an output machine such as a printer is determined by UCR (Under Color Removal) of achromatic component and chromatic component, and a technique called Kueppers Technique is proposed. This technique is the first proposed technique as color conversion processing of the HiFi color and has an easy algorithm and thus is widely utilized.\nConsidering application of Kueppers Technique to color conversion in conformity with ICC, it is possible to adopt a method of converting a device-independent color signal like L*a*b* color space into an RGB color signal by any color conversion unit and performing UCR processing of achromatic component and chromatic component for the RGB signal. However, as with the usual UCR processing, a colorimetric match between the RGB color signal of an input color signal and the YMCKRGB color signal of an image record signal is not considered and therefore the technique cannot realize calorimetric color reproduction and color conversion accuracy is poor; this is a problem. Further, as with the usual UCR processing, the technique involves a problem of incapability of making full use of the color gamut of an output machine and a problem of incapability of making full use of the color gamut widened in the HiFi color.\nIn contrast, the following technique called a division method is proposed: The color gamut of the HiFi color consisting of five to seven colors is divided into color gamut divisions each involving three or four colors in combination, and color conversion from the colorimetric value is performed in each color gamut division according to a similar technique to that with a usual three-color or four-color printer. For example, JP-A-2000-32284 proposes a technique of using color gamut divisions each involving K and other two colors close to K in hue in combination and determining coefficients in a direct look-up table (DLUT) for color conversion processing for realizing colorimetric color reproduction. For example, JP-A-2001-136401 proposes a technique of using color gamut divisions each involving K and other three colors in combination and determining DLUT coefficients for color conversion processing for realizing colorimetric color reproduction. Further, for example, U.S. Pat. No. 5,563,724 proposes a method of forming k (k<n) color gamut divisions each consisting of m color inks in an image formation apparatus of n (n>4) color inks and creating a printer model for each color gamut division, thereby converting from CIE color space into n-color inks color space so that a calorimetric match is found in the color gamut division and the color difference becomes the minimum outside the color gamut division. Further, for example, U.S. Pat. No. 5,892,891 describes a color conversion method from the colorimetric value in five-colors process of YMCKO; it proposes a method of executing color conversion of the colors in the YMCK color gamut from a colorimetric value. A color in the YMCK color gamut is converted from L*a*b* color space to YMCK. A color outside the YMCK color gamut is color-converted using a color gamut division to which a solid color is added (YOMK color gamut).\nEach of the division methods disclosed in the documents has the advantage that the color conversion accuracy is high because colorimetric color reproduction in each color gamut division is insured. Although the computation amount is large as compared with Kueppers Technique, the division method is rarely used directly to perform color conversion; usually, it is used only for determining the lattice point parameter of DLUT for executing color conversion from L*a*b* color space to color signal of five to seven colors of an output machine, and therefore the large computation amount does not introduce any problem.\nThus, the division method has the advantage that the color conversion accuracy is high because colorimetric color reproduction is insured. However, the division method involves an intrinsic problem such that the color separation values become discontinuous on the boundary between the color gamut divisions and a pseudo contour occurs, as described in JP-A-2001-136401. Then, in the method disclosed in JP-A-2001-136401, occurrence of pseudo contours is suppressed by performing smoothing processing of the lattice point parameters of DLUT determined by the division method. However, a new problem is caused, that is, the smoothing processing worsens the color conversion accuracy. That is, in the division method, prevention of a pseudo contour and realizing high color conversion accuracy cannot be made compatible; this is a problem of the division method.\nFurther, a problem occurs in addition to occurrence of the pseudo contour in the boundary depending on the division method into color gamut divisions. In the color gamut division methods disclosed in JP-A-2000-32284 and U.S. Pat. No. 5,563,724, each color gamut division is made up of two or three adjacent colors containing K. In the division methods, the gray axis becomes a single color of K and thus gray cannot be reproduced in three colors of YMC; this is a problem. Usually, gray reproduced in the K single color is poorer in graininess than gray reproduced in the three colors of YMC and therefore a problem of poor graininess in the proximity of the gray axis occurs in the division method. In the division method, the place where the maximum density is generated becomes a K signal color solid part and the maximum density is lower than that in a YMCK four-color solid part in the YMCK four-color process. Full use of the color gamut in the high-density part that should be able to be reproduced essentially cannot be made; this is a problem.\nIn contrast, in the methods disclosed in JP-A-2001-136401 and U.S. Pat. No. 5,892,891, one of the color gamut divisions is formed of the YMCK color gamut, whereby gray can be reproduced in three colors of YMC and the place where the maximum density is generated can be made a YMCK four-color solid part, so that the graininess in the proximity of the gray axis can be improved and full use of the color gamut of the high-density part can be made. However, in the method, discontinuousness of the color separation values in the boundary between color gamut divisions containing RGB and the YMCK color gamut is very large as compared with that in the division methods disclosed in JP-A-2000-32284 and U.S. Pat. No. 5,563,724. There is a problem of easy occurrence of a pseudo contour."}
{"text": "Field of the Invention\nThe present invention relates generally to structural elements for a building, and more particularly to a device, method, and system for framing openings in building walls.\nDescription of the Related Art\nBuildings, such as houses and other structures, are commonly constructed of frame-constructed walls that separate the inside of the building from the outside, and that separate rooms from one another, as well as that support a roof and possibly upper floors of the building. The frame-constructed walls are formed by so-called stick members, which include vertical studs, horizontal wall plates, and horizontal load-bearing headers or lintels over portals, doors, windows and fireplaces. The stick members for a stick-built wall may be of wood, steel, or other materials. The frame-constructed walls serve as a base for wall coverings, such as wallboard, siding, wall sheathing, etc.\nThe walls typically have openings for doors and windows, for example, that are formed in the wall and provided with framing that forms a rough opening for the window or door and that supports the window or door in the wall. In one example, a framed rough window opening has a vertical king stud on each side of the opening extending from the bottom or sole plate of the wall to the top plate of the wall structure. Vertical jack studs are provided below the window opening and vertical cripple studs are provided in the wall above the window opening. Horizontal members of the rough window opening include a sill at the bottom of the window opening and a header or lintel extending horizontally across the top of the opening. Vertical trimmer studs are provided on the sides of the opening extending from the sill to the header or from the sole plate to the header. Doors may be framed in the same way but without the use of a sill and jack studs that form the wall below a window.\nEvery structural opening of a standard stick-built structure requires at least two king studs, two jack studs, and one header. Typically, these three elements are constructed using solid virgin lumber, the construction of which requires multiple calculations and cuts by the contractor, resulting in pieces of scrap wood and using valuable time. These calculations are based off of the expert knowledge that the professional construction worker has gained. Following this process, the three elements are commonly nailed together to form the opening. In the majority of new construction, walls are framed flat on the floor and then stood up and set in place. With this form of framing, doors and windows are laid out first.\nThe process of remodeling, although different, also places a major focus on the importance of structural openings. With the top plates and sole plates already installed, studs are cut and installed vertically, to create new openings.\nIn both cases, the king studs are installed first. The king stud is the piece of lumber that runs from the bottom plate or sole plate of the wall (a horizontal piece of wood usually 1.5 inches thick that acts as the base of the wall to which vertical pieces may be attached) to the top plate of the wall (a horizontal piece that runs parallel to the sole plate that forms the top of the wall). In most states, building codes require that the top plate is to be doubled, making it a total of 3 inches thick. The purpose of the king studs is to uphold the structural integrity of the stud wall and to hold the header in place laterally. King studs can be purchased in pre-cut lengths for standard ceiling heights. For example, for a ceiling height of 8 feet (96 inches), a pre-cut king stud is 92⅝ inches in length. This accounts for a 1.5 inch thick sole plate and a 3 inch top plate. This same calculation is applied to ceilings that are 9 feet and 10 feet high.\nSecond components of a structural opening are the jack studs (also called trimmers). Jack studs are cut and nailed at the inside surfaces of the king studs relative to the wall opening. The jack studs also rest on the sole plate of the wall. The jack studs are responsible for carrying the header and for distributing the bearing weight that the header carries, down to the base plate, ensuring structural stability. The jack studs are typically cut on the job site, and are measured depending on the door, window, or opening height. For example, for a 6 feet, 8 inch (80 inch) tall door, the jack studs are cut to 80.5 inches and installed on the sole plate. It is common knowledge to construction professionals that in order to fit a 6 feet 8 inch door, the opening in the wall must have a height of roughly 82 inches. The same calculation applies to doors that are 7 feet tall and 8 feet tall. Both jack studs and king studs of a typical building structure are formed of flat pieces of conventional lumber.\nThe final component of the three components of a structural opening, according to certain embodiments, is the header. The header can be described simply as the structural piece that spans the opening formed by the jack and king studs. The header is placed on top of the jack studs. Whatever weight the header bears is therefore transferred to the jack studs and down to the sole plates. The header is typically placed directly under the top plate, between the two king studs, and bearing on the two jack studs. The amount of weight that the header bears can vary as it is responsible for supporting the weight of the ceiling, as well as any floors above. The load can often determine the height, width, and composition of the header. Headers are normally built using solid conventional lumber. However, headers can also be of engineered materials, sometimes including sandwich structures of steel or plywood depending on the structural strength required. If the only bearing responsibility for the header is to support sheet rock or other wall portion above the header, a hollow header may be used to reduce cost and save material. Headers are easily manipulated and can also be insulated in order to increase the R-value (insulating value) of that area of the building structure. The R-value is the ability of a given space to resist heat flow. A high R-value means less escaping heat or influx of heat when an air conditioner is in use, which can lead to increased savings on energy bills.\nThe typical necessary dimensions (especially width) of the header are also common knowledge to house construction professionals, for example. The two 1.5 inch jack studs must be accounted for. For a 24 inch wide door, which requires a rough 26 inch wide opening, the header must be cut to 29 inches. When the opening is much larger, it is common for two jack studs to be used on either side of the opening. As a result, the length of the header must account for an added 3 inches on each side, for a total of 6 inches of additional length."}
{"text": "1. Field of the Invention\nThe present invention relates to an imaging apparatus, imaging optical unit, and imaging system such as a video camera, television camera, and digital still camera.\n2. Description of the Related Art\nWhen a subject image is picked up through an imaging optical system lens, the peripheral brightness of the picked-up image is lower than the central brightness. This phenomenon is mainly caused by the cosine biquadratic rule and vignetting.\nIn order to correct such a drop in peripheral light quantity, for example, as proposed in Japanese Unexamined Patent Publication No. 165023 of 1994, it is considered that a ratio of the peripheral light intensity to the central light intensity on the light receiving surface of an image pickup device is determined based on the positions of a zoom lens and a focus lens that comprise an imaging apparatus and a stop value, the reciprocal of the determined ratio is outputted as a correction coefficient, and an image signal is multiplied by the correction coefficient, whereby the light quantity distribution on the image pickup surface is electrically corrected.\nHowever, the method for correcting a drop in peripheral light quantity proposed in the above-mentioned publication is not sufficient for the following reason.\nThat is, in the image pickup device, the output value is not always in proportion to the incident light quantity, and the output value has a peculiar gamma property as shown in FIG. 8. Therefore, even if the peripheral incident light quantity is 10% lower than the central incident light quantity, it is not proper that 10% correction of the entirety of output from the image pickup device is corrected. That is, in many cases, the actual brightness distribution in the shooting range cannot be accurately reproduced.\nAnother method is also considered in which gamma correction is made for outputs from the image pickup device so that the gamma property becomes a straight line, and then the light quantity correction is made. However, the desired gamma property of the imaging apparatus is not a straight line, so that in this case it is necessary to make gamma correction again, resulting in an increase in noise.\nFurthermore, the image pickup device has a characteristic by which the light receiving ratio changes depending on the incidence angle of an incident light flux. That is, the device has light receiving ratio distribution in accordance with the position of the exit pupil, and this also causes the light quantity to lower at the peripheral portion at which the incident light flux greatly tilts with respect to the light receiving surface of the image pickup device. Specifically, this light receiving ratio distribution greatly influences an image pickup device provided with micro lenses at the light receiving surface."}
{"text": "Sputter targets are used in physical vapor deposition (PVD) techniques in order to provide thin films on complex integrated circuits. In these processes, it is important that the sputter target provide film uniformity and minimal particle generation. Electrical resistivity of the coated thin film must also be acceptable.\nIn many prior art techniques for making sputter targets, target densities are low with concomitant increased porosity serving as a possible contaminant site. In order to obtain suitable high density and low porosity, high temperature manufacturing processes are employed that, in turn, sometimes lead to undesirable grain growth in the target.\nOften high oxygen content targets result from the conventional processes. Films coated by such targets exhibit undesirable electric resistivity. Thus, many of the prior art sputter target manufacturing methods result in unacceptable target and resulting film characteristics such as contamination, low density, and undesirable electrical properties."}
{"text": "Traditional techniques for designing and simulating computer-generated environments, such as virtual worlds in massively multiplayer online (“MMO”) games, do not readily facilitate the efficient development and production of new or modified virtual worlds. Moreover, conventional game platforms for generating virtual worlds further impede the development and production of virtual worlds, especially with respect to massively multiplayer online games. Generally, conventional game platforms use a set of game servers to create a virtual world, with each additional set of games servers being used to create a copy of the virtual world. Such copies are typically referred to as “shards.” Conventional game platforms routinely divide the geography of each virtual world into partitions, and allocate one game server to generate each single partition. A game server that is configured to generate a partition of a virtual world is commonly referred to as an “area server.” Further, some conventional game servers are monolithic, which means that all the game processes for generating a partition for a virtual world are instituted in a single process and/or server. While functional, there are a variety of drawbacks associated with conventional game platforms, as well as the traditional approaches to developing computer-generated worlds.\nA drawback for some conventional game platforms is that they allocate one game server for a partition of the virtual world. As such, the game server and its functionality are usually linked directly to the partition and its activities. The performance of a geographically-linked game server, therefore, is a function of the number of players at that partition as well as the kinds of activities being performed there. Further, conventional game platforms that use monolithic game servers are somewhat limited by their inability to adjust their hardware capabilities to meet increases in demand. So, as more players move to a specific partition, the game server that supports that partition becomes increasingly burdened. This, in turn, limits the number of players that can interact with the corresponding virtual world partition. As a remedy, some massively multiplayer online games operate to migrate players off game servers that are heavily populated to game servers that are relatively less populated. This may significantly impact player experience.\nYet another drawback to traditional game platforms is that each game server usually has a separate process space for generating a virtual world partition than the other game servers in the same virtual world. It is not uncommon that players interact across a geographic boundary that demarcates two virtual world portions. As such, difficulties arise due to partition boundaries and conventional game platforms generally require computational resources to try to manage cross-partition interactions. For example, a game server associated with one player in one partition generally requires some sort of tracking functionality to track the interactions of that player with other players in other virtual world partitions. As another example, consider that computational resources are generally required to maintain players as proxies in different process spaces. Traditionally, implementing players as proxies requires converting one representation of player-specific information for one virtual world partition, such as a character identifier, into another representation for use with a separate process space for another virtual world partition. Further, game servers usually communicate among themselves about player interactions using, for example, remote procedure calls, which increases the complexity of communications as well as latency for cross-partition interactions.\nA further drawback to conventional game platforms is that players (or proxies of players) that interact across different virtual world partitions can readily disrupt transactions to unjustly enrich themselves. To illustrate, consider that a first player at one virtual world partition offers to exchange a weapon for gold possessed by a second player at another virtual world partition. Next, consider that the first player receives the gold, but before surrendering the weapon, the first player breaks the connection between the game servers at the different virtual world partitions. As a result, the first player will receive the gold without exchanging the weapon.\nOne drawback to using monolithic game servers is that if one of the game subprocesses, such as a physics process, of a monolithic game server falters or shuts down, then the entire game server in which the game subprocess resides is also susceptible to faltering or shutting down.\nThere are other drawbacks to using current client-server architectures to design and simulate computer-generated environments for massively multiplayer online games. Consider that along with the development activities required to create new content for new or modified virtual worlds, game designers and developers also regularly modify executable instructions to accommodate the new content and its data. Modifying executable instructions usually increases the development time, especially during quality assurance and reliability testing to weed out bugs and other problems. Further, massively multiplayer online games commonly require the modified executable code to be downloaded to a client prior to interacting with the new content. For instance, it is not uncommon for massively multiplayer online game platforms to implement large data downloads (e.g., 250 MB), which can include the modified executable code.\nIn view of the foregoing, it would be desirable to provide a system, a computer-readable medium, a method and an apparatus for generating synthetic environments in a distributed architecture to, among other things, introduce dynamic content into the synthetic environments, and to further provide for efficient development and production of new or modified synthetic environments."}
{"text": "Organic pollutants from diverse agricultural, municipal and industrial facilities, and waste sites partition preferentially into the soil, water or air phases and spread rapidly throughout the environment. Many of these materials, even in small concentrations, adversely affect life forms, and create serious environmental threats. Toxicity, carcinogenicity, and mutagenicity are the most critical biological properties of a potential pollutant, a very large number of organics of which have been identified by the U.S. Environmental Protection Agency as particularly threatening. For example, the widespread use over the last few decades of herbicides, pesticides and related high risk chemicals, e.g., organochlorines, polychlorobiphenyls (PCB's) and chlorinated phenols, have resulted in serious environmental metal problems. In application Ser. No. 519,793 filed May 7, 1990 by David D. Friday and Ralph J. Portier, there is disclosed a continuous flow, immobilized cell reactor, and bioprocess, for the detoxification and degradation of volatile toxic organic compounds. The reactor is closed, and provided with biocatalysts constituted of specific adapted microbial strains immobilized and attached to an inert porous packing, or carrier. A contaminated groundwater, industrial or municipal waste, which is to be treated is diluted sufficiently to achieve biologically acceptable toxicant concentrations, nutrients are added, and the pH and temperature are adjusted. The contaminated liquid is introduced as an influent to the closed reactor which is partitioned into two sections, or compartments. Air is sparged into the influent to the first compartment to mix with and oxygenate the influent with minimal stripping out of the toxic organic compounds. The second section, or compartment, is packed with the biocatalyst. The oxygenated liquid influent is passed through the second compartment substantially in plug flow, the biocatalyst biodegrading and chemically changing the toxic component, thereby detoxifying the influent. Non toxic gases, and excess air from the first compartment, if any, are removed through a condenser located in the overhead of the reactor. Liquids are recondensed back to the aqueous phase via the condenser.\nIn the reactor described in the application, supra, the air in the form of a high velocity jet is sparged into the first compartment of the reactor and combined with the conditioned liquid influent at high gas/liquid shear conditions, under pressure, to create very fine bubbles. The mixed phases of air and liquid are flowed through a bed packed with a solid inert packing, e.g., glass beads, then through an open tubular column to a packed bed of biocatalyst. Good mixing, with minimum stripping of the toxic organic compounds from the liquid, is obtained. The pressure increases the oxygen driving force, and the high gas shear provided by the method of contact between the gas and liquid, and solid inert packing, under pressure minimizes bubble diameter and increases the interfacial transfer area between the phases. By small bubble formation and good mixing with good air/liquid contact at elevated pressure, the volume of air that is introduced into the reactor is minimized. Consequently, a lesser amount of the volatile organic compounds are stripped from the liquid and a greater concentration of the volatile organic compounds are contacted, with the biocatalyst and mineralized to detoxify the influent stream. Whereas the apparatus described in this application has proved admirably suitable for admixing the liquid and oxygen influent phases introduced into the reactor, with minimal stripping of the volatile organics by the gas, there nonetheless remains need for alternate bubble generation devices, as well as a need for improved bubble generation devices. There also remains a need for reactors of improved design for carrying out these types of reactions due principally to the volatility of many of the chemical toxicants targeted for detoxification biodegration."}
{"text": "As the demand for high speed broadband networking over wireless communication links increases, so too does the demand for different types of networks that can accommodate high speed wireless networking. For example, the deployment of IEEE 802.11 wireless networks in homes and business to create Internet access “hot spots” has become prevalent in today's society. However, these IEEE 802.11-based networks are limited in bandwidth as well as distance. For example, maximum typical throughput from a user device to a wireless access point is 54 MB/sec. at a range of only a hundred meters or so. In contrast, while wireless range can be extended through other technologies such as cellular technology, data throughput using current cellular technologies is limited to a few MB/sec. Put simply, as the distance from the base station increase, the need for higher transmission power increases and the maximum data rate typically decreases. As a result, there is a need to support high speed wireless connectivity beyond a short distance such as within a home or office.\nAs a result of the demand for longer range wireless networking, the IEEE 802.16 standard was developed. The IEEE 802.16 standard is often referred to as WiMAX or less commonly as WirelessMAN or the Air Interface Standard. This standard provides a specification for fixed broadband wireless metropolitan access networks (“MAN”s) that use a point-to-multipoint architecture. Such communications can be implemented, for example, using orthogonal frequency division multiplexing (“OFDM”) communication. OFDM communication uses a spread spectrum technique distributes the data over a large number of carriers that are spaced apart at precise frequencies. This spacing provides the “orthogonality” that prevents the demodulators from seeing frequencies other than their own.\nThe 802.16 standard supports high bit rates in both uploading to and downloading from a base station up to a distance of 30 miles to handle such services as VoIP, IP connectivity and other voice and data formats. Expected data throughput for a typical WiMAX network is 45 MBits/sec. per channel. The 802.16e standard defines a media access control (“MAC”) layer that supports multiple physical layer specifications customized for the frequency band of use and their associated regulations. However, the 802.16e standard does not provide support for multi-hop networks.\n802.16 networks, such as 802.16j networks, can be deployed as multi-hop networks from the subscriber equipment to the carrier base station. In other words, in multi-hop networks, the subscriber device can communicate with the base station directly or through an intermediate device.\nThe complexity involved in supporting multi-hop networks in a robust manner necessarily involves sophisticated MAC control layer protocols. Such protocols do not exist. For example, as noted above, the IEEE 802.16e standard does not support multi-hop networks. The IEEE 802.16j standard for supporting multi-hop networks has been proposed, but the standard currently makes no provision for efficient use of MAC layer resources. As such, MAC protocol data units (“PDUs”) in a multi-hop environment are not arranged to minimize overhead or provide efficient means for relaying control information. For example, current methods do not allow MAC PDUs for multiple connections associated with different users to be grouped into a single relay packet. As another example, current methods do not define how route control or quality of service (“QoS”) control within a multi-hop communication network can be supported in centralized or decentralized manners.\nIt is therefore desirable to have method and system that provides MAC PDU arrangements that allow efficient use of MAC layer resources in supporting wireless multi-hop relay networks, including but not limited to those operating in accordance with the IEEE 802.16 standards."}
{"text": "1. Field of the Invention\nThis invention relates to locating RFI sources on power distribution and transmission systems.\n2. Description of the Prior Art\nSparks on/in power system equipment are a significant source of radio frequency interference (RFI). When these RFI sources are of sufficient strength they can disrupt radio and TV reception and all types of radio communications. This leads to complaints to the operator of the power system and the source must be found and corrected. Often the RFI source is a faulty device or appliance in the home or office of the complaining party. However a significant percentage of RFI sources are caused by faulty devices or equipment on the power distribution/transmission system itself. Generally these sources are found on the overhead cables and associated equipment that distribute commercial electrical power.\nAfter receiving a complaint the power company personnel typically locate the pole or structure from which the RFI is originating with special radio direction finding (RDF) equipment. Then they use very short range radio or ultrasonic detectors to pinpoint the source. For many years attempts have been made to design ultrasonic detectors with sufficient sensitivity, directivity and reliability to locate the spark RFI source from the ground, thus eliminating the need for a bucket truck and line crew to locate the fault. The limiting factor has been the achievable sensitivity and signal to noise ratio (S/N).\nThe present invention is a pinpointer, locator or detector of RFI event sources, with increased discrimination of power system RFI sources from random, ambient and other sources. This improved discrimination provides benefits in reducing the time and cost required to locate an RFI source on a power system. This and other improvements are achieved by taking into account the specific ultrasonic characteristics, or signatures, of power system line RFI sources so that they can be discriminated from background noise.\nAs in the prior art, some embodiments of the invention use a parabolic reflector to focus ultrasonic wave energy on a microphone mounted at the reflector focal point, although alternatively other types of ultrasonic sensors (such as direct contacting stethoscope probes) can be used. Unlike prior art, which typically treated the fundamental technical problem as one of detecting certain ultrasonic frequencies emitted by spark sources and mixing them down to provide audio indication to the operator, the concept of the invention evolved from understanding that the ultrasonic energy coming from the spark is essentially a series of shock waves in the atmosphere, rather than any particular ultrasonic frequency, and that these shock waves arrive at the microphone in predictable patterns directly related to the frequency of the AC power system and to the type of interference source from which they emanate.\nThe shock waves are converted to an alternating electrical voltage by a piezo ceramic transducer that resonates at an ultrasonic frequency when stimulated by the atmospheric shock wave. This signal is then amplified and filtered to remove out of band noise such as internal electronic noise and external radio noise picked up by the circuits. The signal then goes to a logarithmic amplifier/detector circuit, the output of which is a voltage proportional to the log base 10 of the AC input signal level. The time constants in the detector are adjusted so that it amplitude demodulates the ultrasonic frequency, thereby providing a signal with bandwidth from DC up to an upper cutoff frequency chosen for optimum system performance. The AC output of the amplifier/detector is sent to the AC amplifier and the bandpass filters, the outputs of which are converted to proportional DC voltages by active rectifiers. (The filter pass bands are designed for the frequency of the power system that the instrument will be used on, generally either 60 HZ or 50 HZ. If the instrument is to be used on more than one type of system the filter pass bands can be designed to be switchable.) The DC voltages from the bandpass-rectifier paths are then compared to each other and the DC output of the AC amplifier-rectifier path by a bank of comparator circuits. The resulting high/low comparator outputs are combined and screened logically so that when a predefined condition is met, based on prior analysis of signals collected from power line RFI sources, an indicator light is turned on to notify the operator. At the same time the unfiltered outputs of the bandpass-rectifier circuits, now providing higher harmonics of the power line frequency in the audible range, are fed to an audio amplifier and speaker. The sound amplitude and timbre from the speaker advise the operator, in a relatively gross way, of what is coming in through the microphone. At the same time the output of one or more of the bandpass-rectifiers (depending on mode of operation) is used to drive a signal strength meter which gives the operator a more precise indication of the strength of the power line component of the received signal.\nThe meter and light indicators are attached to the parabolic reflector in such a way that they are easily visible to the operator while sighting to point the highly directional parabolic microphone unit.\nWhen practicing the invention the operator will usually have identified a particular pole or structure from which the RFI is emanating using RDF equipment. Standing at ground level, he/she will then point the ultrasonic parabolic microphone with one piece integrated electronics, meter, indicator and speaker, at various suspect devices on the pole until a spark source is identified by the speaker and meter and confirmed by the indicator light. The operator will then scan the pole from various angles to confirm which of various closely spaced items is the faulty component. Typically the operator will have to work with high levels of ambient ultrasonic noise, such as from traffic, wind, rustling leaves, etc. Since these noise sources will not have the. specific spectral characteristics of power line sources the indicator will not illuminate (or sound in the case of an audio indicator). Also, because of the response characteristics of the logarithmic amplifier/detector, the speaker and meter will both tend to quiet in the presence of loud out-of-band noise, further enhancing ease of operation.\nIt is anticipated that the invention will be used mostly for locating spark RFI sources, but the same fundamentals apply to corona RFI. In fact, since there are predictable differences in the signals emitted from sparks and corona, a variation of the invention can be used to distinguish between them. For example, a transmission tower may have a spark RFI source which is interfering with local UHF reception and at the same time, because of the high voltage, a corona RFI is present which is chronic but not a source of complaint. The operator has located the tower with RDF equipment which is not precise enough to locate the specific faulty device causing the spark. Using the invention he can select a logic mode to xe2x80x9cseexe2x80x9d only spark type sources. This was not possible with the prior art.\nAdditional features and advantages will be made apparent from the following detailed description of the illustrated embodiment which proceeds with reference to the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus, a circuit and a method wherein an image is transformed by means of texture mapping.\n2. Description of the Related Art\nIn the field of computer graphics, as a technique of transforming an image in a virtual three-dimensional space, a technique called texture mapping is used in which a model that is an object shape and is configured with a combination of triangles called a polygon is made and then an image is attached to the model. Here, in the case where the model is positioned farther than a camera viewpoint in a virtual three-dimensional space, the model is reduced in size.\nOn this occasion, texture mapping with less aliasing is executed by preparing in advance texture data in which the same image is made to have lower resolution in stages (discretely), and texture data whose reduction ratio is close to that of the model is attached to the model. Typically, an anti-aliasing method according to this technique is called the mipmap method, and texture data lowered in resolution in phases are called mipmap images.\nIf the reduction ratio in the original state in which an image has not been reduced in size is defined as 1.0 and the reduction ratio when an image has been reduced in size is defined as 0.0 to 0.999 . . . , in the mipmap method, image data with a finite number of reduction ratios reduced by ½n (one divided by the n-th power of two), such as 1.0, 0.5, 0.25, 0.125, 0.0625 . . . , are prepared as texture data in a memory. After that, if the reduction ratio of a model is, for example, 0.75, the mipmap data whose reduction ratio is 1.0 and the mipmap data whose reduction ratio is 0.5 are read out from the memory, and further, by linearly interpolating these mipmap data by means of the weighting addition of one half each, texture data whose reduction ratio is 0.75 are calculated (refer to Patent Literature 1, for example).\n[Patent Literature 1] Japanese Patent Publication No. 2002-83316 (Paragraph No. 0004)\nHowever, there have been such problems as the following (1) to (3) in executing anti-aliasing in a typical manner by means of the mipmap method.\n(1) If the reduction ratio of a model is of other reduction ratios than mipmap data, anti-aliasing may not be appropriately executed, because two mipmap data sets are linearly interpolated as described above. Thus, there occurs aliasing or a blur in an output image.\n(2) Since it is necessary to prepare texture data of input image, which have been reduced in size by ½n, processing time becomes longer and the scale of a circuit is enlarged.\n(3) Since approximately twice the amount of input image is required as a texture data amount, such as an input image+(half the input image)+(a quarter the input image)+(an eighth the input image)+ . . . , a memory having large capacity is required."}
{"text": "1. Field of the Invention\nThe present invention relates to a gas sensor of Severinghaus type using a hydrogen ion sensitive FET transducer of the structure of an oblong gate-insulated field effect transistor having its gate part at the front end and its electrode part at the other end.\n2. Description of the Prior Art\nThe measurement of concentration of gas such as carbon dioxide or ammonia gas is, of course, important in industrial applications, but recently in the medical field, the measurement of partial pressure of gas in living bodies begins to be taken seriously. For example, in medicine, continuous measurements of gas partial pressure in blood of anesthetized patients, those with advanced diseases or those in the convalescent stage are lending themselves to discovery of emergent situations. For such purposes, a very small gas sensor with less than 2 mm diameter which may be inserted in any blood vessel or muscle tissue is required.\nFor the aforementioned purpose, heretofore in use has been a gas sensor of the Severinghaus type using a very small glass electrode. Minuaturization of the glass electrode, however, is known to involve following problems:\n(a) Because the resistance of the glass film is approx. 10 M.OMEGA., an amplifier with a high input impedance is necessary. PA1 (b) The glass film, being thin, has low mechanical strength. PA1 (c) Because of the electrode area being small, the resistance of the glass film is large. PA1 The measuring instruments should be large and complex and the electrodes themselves are brittle and tend to break. Therefore, especially as a sensor to be inserted into tissues of living bodies for measuring gas partial pressure in living bodies, this instrument has posed a problem in practical applications. PA1 (a) Because of the use of the solid electrode, sensor is inflexible. PA1 (b) Its electric resistance is increased due to the miniaturization. PA1 (1) The photosensitivity of pH-ISFET exists not only on its gate part but on whole of the Si substrate including its electrode part. Accordingly, it is difficult to totally eliminate the photosensitivity by coating only the gate part with a black polymer, as above described. PA1 (2) In the Severinghaus gas sensor, the partial pressure of gas is measured by utilizing the minute change of pH of the gas absorbing liquid. Accordingly, it is necessary to avoid changes of pH of the gas absorbing liquid due to factors other than the gas partial pressure. However most of such black pigments, being chemically unstable, will become acidic due to autoxidation, etc., during a long period of storage, often inviting changes in characteristics as a gas sensor. Accordingly, using the pH-ISFET with its gate part coated with a hydrophilic polymer colored black for allaying the photosensitivity of the pH-ISFET is problematical in practical applications as a gas sensor.\nOn the other hand, a carbon dioxide gas sensor using a solid pH electrode of metal oxide in place of the glass electrode is disclosed in U.S. Pat. No. 3,719,576, etc. This sensor is smaller and slenderer than that using a glass electrode and is, therefore, suitable as a sensor to be inserted in living tissues, but has the following disadvantages:\nOn this ground, the miniaturization had its limit.\nThis problem has been solved by making use of a hydrogen ion sensitive FET transducer (hereinafter referred to as pH-ISFET) having the gate-insulated field effect transistor structure described in U.S. Pat. No. 4,218,298 in place of the glass electrode of the solid electrode. Such a gas sensor using the pH-ISFET is, as described in U.S. Pat. No. 4,409,980, a gas sensor composed of a pH-ISFET and a reference electrode deposited on the substrate in proximity to the gate part of the pH-ISFET, an insulator tube housing this pH-ISFET and the reference electrode with lead wires connected thereto, the gate part of the aforementioned pH-ISFET being located at an opening part provided in the insulator tube, and the lead wires extended along the tube, an electrical insulation resin closing the tube by filling the space of the part inside the tube wall housing the lead wire connecting parts, a hydrophilic polymer layer containing electrolytes which undergo change in hydrogen ion concentration, as it absorbs the gas, and which is placed around the gate part of the pH-ISFET and the reference electrode, enveloping both of them, and a gas permeable membrane coating at least whole of the aforementioned polymer layer.\nThe aforementioned gas sensor is preferable as a gas sensor to be inserted in living bodies. However, the pH-ISFET usually has a photo-sensitivity of the order of 10.sup.-5 -10.sup.-6 V/lux. Therefore, when a gas sensor using the pH-ISFET is used for monitoring, etc., during operation, when it receives illumination of several thousands luxes - several tens of thousands luxes, the photosensitivity of the pH-ISFET poses a serious obstacle to its practical use.\nCoating the gate part of the pH-ISFET with a hydrophilic polymer colored black for allaying the photosensitivity of the pH-ISFET is described in U.S. Pat. No. 4, 273,636. However, the aforementioned proposal is to allay the photosensitivity of ion sensor using ISFET and does not relate to a gas sensor. It is in principle practicable to allay the photosensitivity of a gas sensor by utilizing a pH-ISFET with a hydrophilic polymer colored black coated on its gate part. In this instance, the gate part of the pH-ISFET and the reference electrode are enclosed by a hydrophilic polymer containing a gas absorbing liquid colored black and the outside of this polymer is coated with a gas permeable membrane. It has become evident, however, that diminishing the photosensitivity of a gas sensor by this method involves the following problems:\nThe present inventors have found out the improved gas sensor of this invention as a result of investigations carried out for providing a practically useful gas sensor in which the photosensitivity of the gas sensor of the Severinghaus type using pH-ISFET is allayed."}
{"text": "1. Field of the Invention\nThis invention relates to the prevention of microbially induced corrosion, microbial fouling and ordinary electrochemical corrosion by the application of a coating of nickel-63 to the substrate. Specifically, this invention relates to prevention of corrosion of metallic surfaces by microbes.\n2. Description of the Prior Art\nMicrobial corrosion, or microbiologically influenced corrosion (MIC), may be defined as metal loss caused, or accelerated, by microbial action at one or both of the two sites controlling electrochemical corrosion, namely the anode and the cathode. Microbial corrosion has been noted with cast iron, mild steel, stainless steel, aluminum and its alloys, and copper and its alloys. Concrete, wood and various plastics also undergo microbial degradation (by mechanisms that obviously fall outside the above definition for microbial corrosion). The destruction of metallic and non-metallic materials due to microbial action is a serious and costly problem for industrialized societies.\nA number of microbes have been identified with degradative processes, including those listed below:\nCladosporium resinae--a fungus that corrodes aluminum fuel tanks.\nPseudomonas--a type of bacterium capable of reducing ferric iron to ferrous iron.\nThiobacillus ferrooxidans--a type of bacterium that oxidizes ferrous iron to ferric iron.\nThiobaccillus thiooxidans--a type of bacterium that produces sulfuric acid.\nGallionnella--a type of bacterium that can oxidize Fe.sup.+2 to Fe.sup.+3.\nNitrobacter--a type of bacterium which oxidizes NO.sub.2.sup.- to NO.sub.3.sup.-.\nNitrosomas--a type of bacterium which oxidizes NH.sub.3 to NO.sub.2.\nDesulfovibrio, Desulfomonas, Desulfotomaculum--types of anaerobic, sulfate-reducing bacteria.\nThe mechanisms by which these microbes contribute to the corrosion of metals and degradation of non-metals are unique to the specific microbes, and many are not well understood. The most widely studied form of microbial attack involves the anaerobic sulfate-reducing bacteria (SRB).\nThe SRB are found in a variety of natural environments, including many types of soils and sediments, fresh, brackish and salt waters, natural hot springs, oil and gas wells, and sulfur deposits. They tolerate temperatures from &lt;5.degree. to 75.degree. C., a pH range of about 5 to 9.5, and hydrostatic pressures of at least 10.sup.5 kPa. They survive, but do not grow, under aerobic conditions.\nAlthough there is not universal agreement about the mechanism by which SRB cause or accelerate metallic corrosion, it is believed that the bacteria bring about cathodic depolarization of the metal surface by removing hydrogen from cathodic sites in a sulfate-reducing reaction. The equations for this proposed mechanism are as follows:\n(1) anodic reaction 4 FE.fwdarw.4 Fe.sup.+2 +8 e.sup.-\n(2) electrolytic dissociation of water 8 H.sub.2 O.fwdarw.8 H.sup.+ +8 OH.sup.-\n(3) cathodic reaction 8 H.sup.+ +8 e.sup.- .fwdarw.8 H.sub.2 O\n(4) cathodic depolarization by SRB 8 H+SO.sub.4.sup.-2 .fwdarw.S.sup.-2 +4 H.sub.2 O\n(5) corrosion product Fe.sup.+2 +S.sup.-2 .fwdarw.FeS\n(6) corrosion product 3 FE.sup.+2 +6 OH.sup.31 .fwdarw.3 Fe (OH).sub.2\n(7) overall reaction 4 Fe+SO.sub.4.sup.-2 +4 H.sub.2 O.fwdarw.3 Fe (OH).sub.2 +FeS+2 OH.sup.-\nThe prevention of microbially induced corrosion usually involves trying to prevent the occurrence, growth and metabolic activity of the microbes in the vicinity of the metal by the use of chemical reagents (biocides), cathodic protection, or protective coatings (e.g. an epoxy film).\nU.S. Pat. No. 3,497,434 teaches a method for preventing fouling of a metal structure immersed in a marine environment by coating this metal structure with a metal toxic to marine organism. The metals toxic to marine organisms which are described in this patent are cadmium, tin, zinc and aluminum.\nU.S. Pat. No. 4,123,338 describes a method of preventing fouling coupled with corrosion inhibition in water environments by the application of a coating of technetium 99 to the substrate."}
{"text": "Power electronics are widely used in a variety of applications. Power electronic devices are commonly used in circuits to modify the form of electrical energy, for example, from AC to DC, from one voltage level to another, or in some other way. Such devices can operate over a wide range of power levels, from milliwatts in mobile devices to hundreds of megawatts in a high voltage power transmission system. Despite the progress made in power electronics, there is a need in the art for improved electronics systems and methods of operating the same."}
{"text": "1. Field of the Invention\nThe field of invention includes electrical safety apparatus in combination with a container. More particularly, the field involves lighting apparatus, electrical power extensions, and electrical shock protection devices which are self-containerized for portability, safety and field use. Such devices are in compliance with most State and Federal Rules and Regulations including those of the Office of Safety and Health Administration (\"OSHA\").\n2. Definition/Explanation of Terms\nGround Fault Interrupter (GFI). The term Ground Fault Interrupter is a shortened version of the original, longer, more descriptive terminology, Ground Fault Circuit Interrupter (GFCI). Both terms and acronyms are used interchangeably. Such a device is electrically in series with a power source and any given portable tool which a user may plug into a female receptacle incorporated with a GFI.\nThe purpose of the ground fault interrupter is to protect a user by preventing electrical shock from portable tools and power equipment. A GFI device performs its safety function by sensing minute stray currents (on the order of 4 to 6 milli-amps) which may be attempting to reach earth ground through a human's body.\nThese very small initial currents flowing to ground and detected by the GFI provide a near instantaneous circuit breaker which interrupts the primary electrical current from the power source. This detection occurs within approximately 1/30th of a second before higher, more dangerous current levels can build up, thus preventing electrical shock, reducing fire hazards and creating a safer work place.\nReceptacle. An electrical receptacle is one or more female receivers usually of the three prong variety. They are also normally configured in pairs or multiple ganged pairs. GFI circuitry is incorporated internally and is electrically in series with the internal and external connectors in a housing for the receptacle. This configuration thereby creates an integral circuit interrupter and connection apparatus housed along with the receptacle.\nSelf-Containment. A plastic containment bucket assembly is a preferred form for self containing the electrical apparatus of the invention. Such a bucket provides a structural mounting apparatus and a weathertight container for housing for the GFI/receptacle unit of the invention.\nExternal Source of Potential. In electrical terms, potential is defined as the potential to do work. It is the electrical potential above earth ground--usually measured in volts. For purposes of this self-contained invention the potential source, or primary source, is located externally, and the invention includes a safe, easy and OSHA accepted way of connecting to such a primary source.\nGround. Ground is defined electrically as zero potential and is used interchangeably with earth ground. Grounding is the act of connecting a device, wire or object electrically to a portion of a circuit which is then physically in contact with the earth usually through a conductive rod driven into the ground.\nWeathertight Container. While the self-contained electrical apparatus of this invention is not absolutely waterproof in use, it is essentially so with normal care and usage. Thus, mounted electrical apparatus in combination with an easy ingress/egress container meets weathertight standards in that conventional hinged covers, lids and sealable gaskets are provided for the electrical apparatus. Moreover, the container has easy ingress/egress and is readily resealable by an open/close lid for convenient daily use by a user.\n3. Description of the Prior Art\nAs we have progressed in our technological development, we have created more uses for tools which use electricity at temporary sites. With this proliferation, concerns for electrical safety at the job site have dramatically increased. The first major thrust into this safety area was the U.S. Electrical code requirement of a separate, third ground wire for electrically supplied power.\nThe next, and still current addition in the safety area of prevention of electrical shock is that of a Ground Fault Interrupter (GFI). This GFI device has become an industry standard for electrical shock protection when using portable tools. Safety has become so critical that such devices are mandated by OSHA.\nGFI units are separately packaged devices between a primary power source and portable power tools. Although fine for shock protection, the nature of construction sites poses a major electrical hazards and housekeeping problems associated with GFI assemblies. Such assemblies, in the past, were thrown on the floor, tacked or tied to temporary construction members. And, in many instances, were susceptible to misuse and damage. Often, GFI units are dragged through, or left in dangerous areas, such as those where water has pooled, thus endangering the lives and safety of workers.\nTurning now to the prior art, a search has turned up various patents, many of which are only of peripheral relevance to this invention. Such patents include, for example, U.S. Pat. No. 3,872,354 (Nestor et al, 1975). Circuitry such as disclosed in Nestor et al is incorporated in today's commercial GFIs. Nestor discloses one typical technique and circuitry for a primary circuit interruption with fast response.\nU.S. Pat. No. 4,709,980 (Coll et al; 1987) is of limited relevance to the present invention. The Coll invention is specifically engineered for cable splicing of joined, contained, and underground cables. It is not designed for repeated, daily use. Its purpose is to receive a properly spliced underground cable and permanently store the same. It is not relevant in structure or function to the invention.\nU.S. Pat. No. 5,217,298 (Jackson et al., 1993) and U.S. Pat. No. 3,066,217 (McDonald, 1962) are lighting only container devices and are unrelated to the portability and safety features of my invention. They have no relation to OSHA requirements nor safety considerations for portable electrical equipment at job sites.\nU.S. Pat. No. 4,984,685 (Douglas, 1991) suggests the use of a plastic bucket for storing and removing an electrical extension cord and cord light in a specific retrievable manner. While Douglas does show the usage of a bucket as a container for such storage it does not disclose the novel features of a ground fault interrupter and bucket combination, nor the combination of the bucket as both a mounting means and weathertight storage. Moreover, Douglas does not teach or suggest a container as a light base for a pole light partially housed in the bucket nor as a weathertight container for other tools which may safely be plugged into a bucket- mounted GFI/receptacle unit."}
{"text": "This invention relates to a wire electrode with a sheath that contains a core of a composite powder core for use primarily in thermal spraying and welding processes, the preparation of the cored wire, and method of using it.\nThe problems of wear and corrosion of metallic surfaces, such as the inside of industrial boilers, may be ameliorated by a judicious application of a protective coating. Such coatings generally require chemical compatibility with the substrate, compatibility in thermal expansion coefficients, low porosity, and good adhesion. When the coating is applied, a coating process must be selected to be compatible with the substrate, its topography, surface curvature, and satisfy constraints imposed by line-of-sight limitations, while maintaining a desired thickness and avoiding thermal distortion of the coating and substrate. Additional detail is found in such sources as R. W. Smith, xe2x80x9cThermal Spray Technology-Equipment And Theoryxe2x80x9d, ASM INTERNATIONAL (1992) and T. J. Mursell and A. J. Sturgeon, xe2x80x9cThermal Spraying of Amorphous and Nanocrystalline Metallic Coatingsxe2x80x9d, WORLD CENTER FOR MATERIALS JOINING TECHNOLOGY, March 1999, which are incorporated herein by reference.\nTraditionally, thermal spray techniques have demonstrated their capability to meet many of these requirements. In such processes, molten particles impact onto a substrate surface to the form a coating.\nThermal spray processes have relatively high coating rates due in part to their high-temperature heat sources. Additionally, such processes involve wide compositional ranges compared to other PVD or CVD coating processes. Coating materials formed thereby effectively serve as a new surface material added to the original surface (substrate).\nVariants of thermal spray technology are capable of processing a wide spectrum of materials, thus making this technology an attractive choice of coating method. Metals from aluminum to tungsten, ceramics, polymers, and advanced composite materials are now being thermally sprayed. Interest in thermal spray processes and equipment has been stimulated by the versatility of the process, the demands of hostile operating environments and materials which need cohesive and adhesive coatings, and the development of new materials and surfacing concepts.\nThermal spray is limited, however, in that it is primarily a direct xe2x80x9cline-of-sightxe2x80x9d process where the incident droplets are deposited only onto surfaces that are aligned with the path followed by the incident droplets.\nA wire-arc spray process is a type of thermal spray process that traditionally utilizes a DC electric arc to directly melt consumable electrode wires. This contrasts with other thermal spray processes which indirectly heat the particles with heated jets of gas. Thus, the thermal efficiency of the wire-arc spray is potentially higher than that of other thermal spray processes. In twin wire-arc spraying, the electric current is carried by two electrically conductive, consumable wires. An electric arc is created between the wire tips across a gap created by the continuous convergence of two wires. A gas jet blows from behind the converging electrodes and transports the molten metal that continuously forms as the wires are melted by the electric arc. High-velocity gas breaks up or atomizes the molten metal into smaller particles in order to create a fine distribution of molten metal droplets. The droplet-infused gas then accelerates the particles away from the electrode tips to a substrate surface, where the molten particles impact to incrementally form a coating upon cooling. Unlike combustion or plasma spraying, the droplets are already molten when they enter the jet zone, and are continuously cooled in transit to the substrate as the particles leave the arc zone. Thus, the optimized wire-arc spray process attempts to shorten the time in transit so that the cooling effects on the molten droplets do not reduce the amount of molten particles that are needed to form continuous coating layers.\nParticles created by most conventional wire-arc spray processes tend to be larger and more irregularly sized than those found in powder-fed thermal spray coating processes. This size irregularity contributes to unwanted deposit porosity, which is common in many wire-arc spray coatings. Droplet atomizing irregularity also depends on uniformity in wire feed rates, the stability of electrical voltage and current, and nonuniformity in the arc gap. The attainment of high-density wire arc coatings is to some extent possible with careful control of wire feed rates, the use of smaller diameter wires, the application of lower currents and wire feed rates, and by experimenting with the atomizing gas-to-metal feed rate.\nThe consumable two wire-arc technique is responsible for the growth of metallizing as a commercial coating method, especially using aluminum and zinc alloys. Smith, supra, p. 23. Despite the wire arc\"\"s high deposit rate capability, it is limited to materials that are conductive and to those materials which are ductile enough to be formed into wire. This reduces the number of materials that are suitable for use in wire-arc processes. However, recent cored wire development has expanded the range of materials to include cermets and xe2x80x9chardfacingxe2x80x9d carbide/oxide with metal matrices. Id. However, such advances still leave the unsolved problems of adhesion to the substrate and cohesion within an often porous coating.\nMechanical, e.g., folded, abutted, or lapped seams of the tube or sheath that contains a core powder found in conventional electrodes present problems when used in welding and thermal spray applications.\nThe longitudinal seams of sheaths formed in conventional processes are closed by the compression stress, which results from reduction operations. Once the forming and sizing are complete, these cored wire electrodes are traditionally wound on spools for feeding into the wire or welding apparatus, or into coils or some other packaging, including drums that contain the wire for stowage, shipping and application. This winding process may relieve the mechanical compression stress that provides the seam\"\"s integrity, causing the seam to open and release or spill the core contents. If the wire has low ductility, it may break. When the wire is being used for welding or thermal spraying, the apparatus includes a wire feed drive system, conduit/liners or guide tubes, welding nozzle/thermal spray front end and a power source. The wire is fed from its container through a set of drive rolls that apply a pinching effect to the wire, which can cause the seam to open. The wire also passes through a flexible conduit that protects and guides it from the spool or package payoffs to the welding or thermal spray gun (head). Any movement of the gun or flexible conduit will also bend the wire, and may cause the seam to open or brittle wire to break.\nOpen seams reduce the strength of the wire and allow the contents to spill. This tends to clog the conduit and reduce the effectiveness and performance of the torch/gun, create amperage drops, cause unpredictable feed and spray rates, allow binding in the contact tips and complete stoppage of the application.\nAdditionally, inconsistent seam integrity Will create problems with the transfer of electrical current across the anode-cathode gap. Electrical current primarily flows around the surface of the wire. In welding and thermal spray processes, the wire is energized at a point as close as possible to the arc to create the least electrical resistance, consistent current flow and to reduce heating of the wire. The current is transferred to or from the wire through electrical contact tips. The contact tips, which are normally copper or brass, have a very close tolerance and limited contact surface. Any interruption in current flow caused by an irregular shape of the wire will be manifest in complications with the weld deposit or thermal sprayed coating.\nFurther, sealed seams, when their integrity is weakened, may permit moisture to penetrate past the sheath/tube and become absorbed in the core material, which causes porosity in coatings, and porosity and hydrogen embrittlement in welds.\nCoatings deposited by wire arc spray techniques generally contain pockets of metal that are thicker and more varied in size than those found in plasma spray coatings deposited by powder combustion and related conventional processes. Large areas of porosity may result from solid, larger particles being trapped, thereby creating voids that have not been filled by subsequent droplets. The conventional wire-arc spray deposit\"\"s coarser microstructure and porosity are often attributed to the large and irregular droplet sizes caused by wire-atomization irregularities. Wire arc spray coatings tend to be rougher than those formed by other thermal spray processes, especially where large diameter wires are used.\nWire-arc spray technology is suited to depositing layers of corrosion- resistant and conductive aluminum, zinc, their alloys, copper, and stainless steels at high rates ( greater than 15 kg/hr [33 lb/hr]).\nMany municipal and civil structures such as bridges, water storage tanks and marine coatings for ship stacks and decks are coated with wire arc spray apparatus.\nAgainst this background, there remains a need to eliminate or effectively reduce surface problems associated with high levels of oxidation, corrosion, and abrasion in industrial environments, including, but not limited to the power generation boiler environment, the paper and pulp industry digester, petro-chemical and chemical manufacturing processing equipment, turbo machinery get and steam turbine) and industrial combustion engines.\nThese and related challenges have led to the development by the inventors of a multiplex powder composite selected from a group of particle sizes consisting of micron (1xc3x9710xe2x88x926 meters) and sub-micron (less than 1xc3x9710xe2x88x926 meters), including nano-scale (1xc3x9710xe2x88x929 meters) particles for a core material in a sheath that serves as an electrode in a wire arc spraying or welding apparatus.\nThe cored wire electrode manufacturing process involves taking a flat metallic strip material, forming it into a xe2x80x9cUxe2x80x9d shape, filling it with the granular or powder core material, then roll forming it into a tube and reduction drawing or rolling it to a suitable size for welding or thermal spraying applications. See, e.g., commonly owned U.S. Pat. No. 5,479,690 which is incorporated herein by reference.\nThe process for manufacturing a sheath or welded seam product requires changes in conventional procedures in order to address issues that are inadequately addressed by traditional methods. The first problem is how to continuously form a tube and weld the seam. The second challenge is how to get particles in the sheath during the forming and seam welding operation without affecting the quality of the weld or the quality of the core material. Good results have been obtained in manufacturing by orienting the sheath in a vertical position using forming dies, guide rollers, pinch rolls for welding and a fill tube that allow deposit of the core material at a point beyond the welding operation. See, e.g., commonly owned U.S. Pat. Nos. 5,346,116; 5,328,086; 5,479,690, which are incorporated herein by reference.\nThe seamless tube or sheath and the composite particles when sprayed provide a more dense coating than is currently available with conventional wire arc thermal spray materials. The added efficiencies, the ability to gain greater reductions, and core compaction enable a cored wire to be made that contains less than 10% voids by volume and a relatively uniform particle distribution. Additionally, chemistries can be balanced based on the sheath and core material, which enhance the versatility and effectiveness of the wire in producing the desired coating or weld.\nThe advantages of a multiplex composite powder core material that is encased in a seamless tube or sheath include:\nElimination of moisture absorption problems within the core material, which permits the use of liquid and near liquid drawing lubricants and cleaning agents and tends to reduce porosity;\nImproved feeding;\nContainment of the core material, which eliminates powder loss and voids in the multiplex composite;\nGreater wire ductility;\nReduced apparatus down time for maintenance;\nMore consistent electrical current transfer and therefore arc stability;\nEnhanced drawing/reduction capabilities;\nBetter compaction of core material, which manifests itself in a more consistent and balanced coating chemistry, density, and particle distribution; and\nSintering of the core material during an annealing step, thereby eliminating the need for a separate sintering operation."}
{"text": "1. Field of the Invention\nThe present invention relates to an information retrieval system for retrieving information from a data base. More particularly it relates to an information retrieval system improved in efficiency by shortening the waiting time at retrieval and one capable of presenting information, if necessary, by processing the retrieval results of various formats obtained through retrieval so that they may be evaluated easily.\n2. Description of the Related Art\nVarious kinds of inventions have been proposed with respect to the information retrieval technique for obtaining desired data by access to data bases dispersed on a network. A conventional technique of this kind is disclosed in, for example, Japanese Patent Publication Laid-Open (Kokai) No. Heisei 2-87274, \"On-Line Information Retrieval System\". The same publication discloses a system for obtaining desired data from a data base on a network, which comprises a communication means for gaining access to a data base on a network, an input means of a retrieval expression, a conversion means for converting the format of the received retrieval expression into one for use in a data base to be retrieved, and a means for storing and editing the data results retrieved by the use of the converted retrieval expression.\nAnother conventional technique of this kind is disclosed in, for example, Japanese Patent Publication Laid-Open (Kokai) No. Heisei 7-56932, \"Information Retrieval System\". The same publication discloses a system for obtaining desired data from a data base on a network, which comprises a for receiving a retrieval request for analyzing it, a for converting the format of the analyzed retrieval request into that one for use in a data base to be retrieved, a means for selecting a proper data base depending on the content of the retrieval request, and a means receiving retrieval results from the selected data base for evaluating them.\nFurther, various means to improve retrieval speed have been proposed in the conventional information retrieval techniques. A conventional technique of this kind is disclosed in, for example, Japanese Patent Publication Laid-Open (Kokai) No. Heisei 7-28836, \"Data Base Retrieval Method and a Data Base System Using This Method\". The same publication discloses a system for sharing the retrieval processing corresponding to a retrieval request among a plurality of retrieval server tasks so as to execute the retrieval processing in parallel, the system comprising a plurality of external storage operable in parallel and a plurality of retrieval server tasks for retrieving data individually from each external storage. In this way, many of the conventional techniques of this kind aim to improve the speed of the retrieval processing itself.\nFurther, swift and accurate evaluation of retrieval results is an important matter in the conventional information retrieval technique. Various techniques for reducing user's labor by devising the output format of retrieval results have been proposed. A conventional technique of this kind is discloses in, for example, Japanese Patent Publication Laid-Open (Kokai) No. Heisei 7-319905, \"Information Retrieval System\". The same publication discloses a system for supporting the operation of a user's focussing on desired data among the retrieval results, after data retrieval, a means for sorting retrieval results based on a predetermined viewpoint, a means for displaying the sort results simultaneously, and a means for displaying the details of some sort results selected from the displayed sort results.\nThe above-mentioned conventional information retrieval techniques, however, are insufficient to improve the retrieval speed and support the evaluation work of retrieval results.\nWhen considering the access to a data base dispersed on a network, in order to improve the retrieval speed, it is necessary to do the retrieval processing in consideration of the difference in the response time per each data base and the difference in the burden of the network. However, these considerations are taken in none of the above mentioned conventional techniques. Even if the above-mentioned conventional system is used in gaining access to a plurality of data bases dispersed on a network, there is no use expecting substantial reduction in the retrieval time, with necessity of further improvement remaining there.\nSince the system disclosed in the patent publication No. 2-87274 sequentially gains access to all the data bases, it can't make up for delays in the response time in the data bases.\nIn the system disclosed in the patent publication No. 7-56932, though data bases accessible depending on the content of a retrieval request are restricted, no consideration is taken in case of an extremely poor response occurring in some data base of the selected data bases to be retrieved. Therefore, it is insufficient to improve the retrieval speed and the waiting time.\nThe system disclosed in the patent publication No. 7-28836 can shorten the time required for retrieval processing itself. However, when retrieving a plurality of data bases, the system can't display the retrieval results until the retrieval result from the last data base is received, because the system displays the retrieval results sorted and integrated. In case of a poor response in a retrieving data base, the system is provided with no means for interrupting the retrieval processing. Therefore, a user must wait until the retrieval in a data base of the poorest response has been completed.\nWhile, in case of access to a plurality of data bases, it is necessary to eliminate redundancy in retrieval results and standardize the display format of the retrieval results, in order to do efficient evaluation work of the retrieval results. The system disclosed in the patent publication No. 7-319905, however, has no consideration as for the overlapping retrieval results and it can't eliminate redundancy caused by the overlap of the retrieval results.\nThe system disclosed in the patent publication No. 2-87274 has no consideration for the difference in display format when searching for overlapping data, although it is provided with a function to eliminate the overlapping data. When judging whether it is the same retrieval result or not, the difference in the display formats of the retrieval results affects the judgement. Further, error compensation for searched overlapping data is not considered in the system.\nThe system disclosed in the patent publication No. 7-56932 can search for overlapping data by comparison among the obtained retrieval results. However, since it only makes a comparison among a plurality of retrieval results, it is difficult to recognize whether it is the same information as for the retrieval result obtained in a different display format."}
{"text": "This invention relates to an image processing method and apparatus thereof, and more particularly, to an image processing method and apparatus which performs a rounding processing.\nWhen a conversion processing is performed on k input data having an arbitrary bit length, if a conventional processing circuit shown in FIG. 8 is used, output data includes some amount having a positive deviation or negative deviation from a mathematical or theoretical value (hereinafter referred to as a \"true value\"). In such errors, there is an error element which can be canceled out when statistically meaningful data are added or data including an error are integrated over a predetermined period of time (hereinafter referred to as a \"symmetric error\"). There is also an error element in which the error cannot be canceled out by the above operation (hereinafter referred to as an \"asymmetric error\").\nIn the conventional processing circuit shown in FIG. 8, when k input data x.sub.1, x.sub.2, . . . , x.sub.k (whose data lengths are n.sub.1, n.sub.2, . . . , n.sub.k bits respectively) are inputted from an input terminal 1010 in parallel, each of the data is respectively multiplied in conversion units (T.sub.1) 1020, (T.sub.2) 1030, . . . , (T.sub.k) 1040 and a rounding processing is performed. The k output data x'.sub.1, x'.sub.2, . . . , x'.sub.k (whose data lengths are m.sub.1, m.sub.2, . . . , m.sub.k bits respectively) obtained from these conversion units are added together by an adder 1050 and a single data whose data length is m' bits is obtained. The rounding processing is performed on this data by a rounding processing circuit 1060 and the result is outputted as a rounded output X whose data length is P bits.\nIn the above processing, when a predetermined number of data having statistical meaning or data over the predetermined period of time are considered, maximum errors .+-..DELTA.x.sub.1, .+-..DELTA.x.sub.2, . . . , .+-..DELTA.x.sub.k are respectively generated in the output data x'.sub.1, x'.sub.2, . . . , x'.sub.k by the rounding processing in the conversion units (T.sub.1, . . . T.sub.k). If an asymmetric error with respect to the true value is .alpha., and a symmetric error is .beta. (where .alpha. and .beta. are both positive or both negative), the maximum errors generated by the rounding processing of the rounding processing circuit 1060 are as follows: ##STR1##\nAccordingly, the maximum error generated in the output X is a sum of errors generated by the rounding processing in the conversion units (T.sub.1 . . . T.sub.k) and errors generated by the rounding processing circuit 1060, and expressed by the formula (2): ##STR2##\nWhen the processing is performed in accordance with the above conventional approach, as indicated by the formula (2), the maximum error having positive deviation from the true value in the output X is not equal to the maximum error having negative deviation. Accordingly, when the output X is considered on the predetermined number of data or the data over the predetermined period of time, the asymmetric error cannot be canceled and remains contrary to the symmetric error. This remainder appears as a noise component in a signal processing and an S/N ratio of the signal processing is deteriorated as an asymmetric error becomes larger.\nFor example, the case where the processing circuit shown in FIG. 8 is adapted in a color printer connected to a host computer, color image data such as RGB signals inputted from the host computer are inputted to the color printer. A color conversion processing is performed on the RGB signals for color image formation in the printer and the RGB signals are converted to data represented in the other color space (e.g. YUV, L*a*b*, YMCBk). The processing circuit shown in FIG. 8 is used for this color conversion processing. However, as described above, an asymmetric error is generated by the rounding processing in the color conversion processing and this remains as noise in the output signal (the color data after the conversion). In other words, generation of the asymmetry error causes deterioration of color reproducibility in the color conversion processing.\nFurthermore, the color conversion processing is repeated numerous times in the color printer, and the asymmetric errors are accumulated, thus creating a problem in the color reproducibility in the color image data processing."}
{"text": "1. Technical Field\nThis application relates to an improved failure sensing device for use in diaphragm pumps.\n2. Description of the Prior Art\nDiaphragm pumps are widely used in the chemical, petrochemical, process and other industries, as noted for example, in U.S. Pat. No. 3,285,182, which refers to such specific uses as in reagent feeding in reactor systems, replenishers and activators in photochemical and electro-chemical systems.\nU.S. Pat. No. 3,666,379, teaches that when handling corrosive fluids, it is desirable to use a diaphragm disk which should be made of a chemically inert polymer such as polytetrafluroethylene (PTFE). The patent indicates that the previous inability to use PTFE can be overcome through the use of the disk referred to in the patent. Unfortunately, attack of the diaphragm by corrosives, fatigue or solvents still leads to the failure of the diaphragm.\nThe attempts to improve the physical design and chemistry of the material of construction of the diaphragm have lead to improvements. However, diaphragms still have a finite life span.\nDevices in accordance with the prior art in the past have been known to exhibit certain shortcomings and problems. As noted in U.S. Pat. No. 3,816,034, mechanically or positively actuated diaphragm pumps suffer from the disadvantage that the diaphragm, due to its positive mechanical connection to a reciprocating drive, is subject to a combination of high shear, bending and tension stresses in the operation of the pumping cycle. Such stresses then lead to shutdowns and replacement of the diaphragm when destruction thereof is imminent or has occurred.\nThe failure of the diaphragm is normally proceeded by the development of minute cracks, tears or hair line fissures which expand until there is a pathway completely through the diaphragm.\nThe problem associated with the replacement of a diaphragm after failure is that failure takes place when the diaphragm stress induced cracks open to the point where there is a hole completely through the diaphragm and consequently, there is a passage for the flow of liquid between the two pump chambers which are normally separated by the diaphragm. Thus, failure can result in heavy losses due to contamination of material streams, exposure of hardware to corrosives, and excessive down time. The fundamental problem is that it is essentially impossible to predict with accuracy the point in time at which the diaphragm will fail. Predictions are thus made on a statistical average basis, which means that some diaphragms will fail before the time period for periodic replacement and some diaphragm will be replaced prematurely, that is, before the major time period of their life span has been expended."}
{"text": "When disasters (e.g., hurricane, storm, tsunami, etc.) strike, many establishments set up collection mechanisms to collect donations to support the victims. These mechanisms may include links on web sites to encourage donations electronically. For example, donations may be done via credit cards or via any other forms of monetary transfers including transfers from accounts associated with payment facilitators (e.g., PayPal Inc. of San Jose, Calif.). Unfortunately, when there are people making donations to the victims, there are also scammers that take advantage of the situations. These scammers may fraudulently claim that they collect the donations for the victims. Since the donation levels typically are at their highest immediately after the disasters struck, it is often difficult to quickly verify the scammers or where the donations go during that time. By the time the scams are detected, large amounts of donations have already been misdirected to the scammers and away from the victims."}
{"text": "Software may stop working correctly for several reasons including “malware” and “software rot.” Malware is a source of problems on computer systems and it comes in a variety of forms with unlimited vectors of attack. Regardless of the form and vector, the end result of malware is that the original software does not work correctly. “Rot” (also known as “software aging”) concerns the condition that after some period of time software stops working correctly. The reasons could be resource consumption, bugs, or transient hardware faults like stray energetic particles striking a chip. Whatever the reason for the rot, the end result is that the software stops working correctly."}
{"text": "Methods and devices are known for operating a vehicle in which a torque of the vehicle is ascertained.\nToday an engine control unit manages the drive train of the vehicle. Not only are the driver's intentions detected and the engine of the vehicle correspondingly controlled, but also torque requirements or limits of other components of the drive train, such as the driving dynamics control or the vehicle transmission, are taken into account. The engine control unit calculates the torque at various locations on the drive train of the vehicle, such as at the output from the engine, the output from the transmission, and/or at a wheel of the vehicle. These torques are measured during an application phase and stored in characteristic maps. The actual torque present is no longer measured during operation of the vehicle."}
{"text": "In a typical wireless network, many devices can enter an area serviced by a wireless controller and communications can be set up between the devices and the controller. In state of the art systems, significant overhead features and functions are required to connect a device to a network. To facilitate an efficient set up between multiple networkable devices, communications must be effectively configured and managed. Thus, a typical wireless network has a communications coordinator or controller such as an access point, a piconet controller (PNC), or a station that configures and manages network communications. After a device connects with the controller, the device can access other networks, such as the Internet, via the controller. A PNC can be defined generally as a controller that shares a physical channel with one or more devices, such as a personal computer (PC) or a personal digital assistant (PDA), where communications between the PNC and device(s) form the network.\nThe Federal Communications Commission (FCC) limits the amount of power that network devices can transmit or emit. Due to the number of networks, crowded airways, requirements to accommodate more devices and the low power requirements, new wireless network standards continue to be developed. Further, it is known that the path loss for systems that transmit at 60 G Hz is very high. Complimentary metallic oxide semiconductors (CMOS) are popular for manufacturing network components due to their low cost and low voltage requirements and a CMOS type power amplifier is inefficient at such frequencies. Generally, most low power 60 gigahertz (GHz) communication systems require directional beam formed systems to achieve acceptable signal to noise ratios (SNR)s and hence, acceptable communication performance.\nAccordingly, there has been a lot of activity to develop low power network communications in the 60 GHz range utilizing directional communications via millimeter waves. An omni-directional transmission or communication is different from a directional transmission, as an omni-directional transmission generally has a single antenna that provides a point source of radiation. A point source has a radiation pattern where signal energy propagates evenly in a spherical manner unless obstructed by an object. In contrast, a directional communication can focus a signal or steer a signal towards a target and a receiver can focus its sensitivity or steer its sensitivity towards a source. In order to accurately steer, transmit and receive beams in the proper direction, a beamforming process or training process can be implemented between a controller and a device. To achieve directional communications several different directional transmission methods are available. One such method utilizes a sectored antenna approach, which switches signals among several predetermined beams. This method can utilize a phased antenna array where transmit and receive beams are formed by changing the phases of the input and output signals of each antenna element. Another approach distributes the transmit power to multiple power amplifiers and, based on the gain of the power amplifiers, the beam can be adaptively steered.\nBeamforming is important for reliable operation of new, state of the art, high frequency low power networks. It can be appreciated that traditional omni-directional transmissions/communication topologies cannot provide reliable low power, high data rate communications at distances of over a few meters. Generally, directional antennas or antenna arrays can provide gains that are much higher than omni-directional antennas (tens of decibels). When a transmitter accurately focuses signal energy in the direction of the desired receiver and a receiver focuses it's receive sensitivity in the direction of the transmitting source, interference from other directions can be mitigated and the beamformed system will create less interference for other systems.\nA directional transmission system can provide improved performance over omni-directional systems due to the increased signal strengths between devices and decreased interference from devices transmitting from directions where the receiver is less sensitive. Higher data rates, on the order of a few Gigabits per second, are possible in a directional transmission mode since the directional link benefits from higher antenna gains. However, these directional systems are typically more complex, slower to set up and more expensive than traditional omni-directional transmission systems. After an association and beam calibration process, efficient data exchange between the device, the controller and other networks, such as the Internet, can occur.\nIt can be appreciated that many network environments, such as offices, office buildings, airports, etc., are becoming congested at network frequencies as many devices enter a network, exit the network and move in relation to the controller of the network. Setting up directional communication and tracking movement of devices in traditional systems requires a relatively long, inefficient association time and set up time for each device. Such continued increase in the number of users for an individual network, the increase in set-up complexity and overhead continue to create significant problems."}
{"text": "In recent years, a flat panel display is developed briskly. Especially, the liquid crystal display device gets a lot of attention from advantages, such as light weight, thin shape, and low power consumption. In an active matrix type liquid crystal display device equipped with a switching element in each pixel, a structure using lateral electric field, such as IPS (In-Plane Switching) mode and FFS (Fringe Field Switching) mode, attracts attention. The liquid crystal display device using the lateral electric field mode is equipped with pixel electrodes and a common electrode formed in an array substrate, respectively. Liquid crystal molecules are switched by the lateral electric field substantially in parallel with the principal surface of the array substrate.\nOn the other hand, another technique is also proposed, in which the liquid crystal molecules are switched using the lateral electric field or an oblique electric field between the pixel electrode formed in the array substrate and the common electrode formed in a counter substrate."}
{"text": "1. Field of the Invention\nThe present invention relates to poinsettia-like ornaments and more particularly pertains to a new decorative ornamental assembly for providing easy to use and colorful decorations about Christmas trees and on building structures.\n2. Description of the Prior Art\nThe use of poinsettia-like ornaments is known in the prior art. More specifically, poinsettia-like ornaments heretofore devised and utilized are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nKnown prior art includes U.S. Pat. No. 1,938,736; U.S. Pat. No. 5,063,485; U.S. Pat. No. 4,171,754; U.S. Pat. No. 3,146,955; U.S. Pat. No. 1,440,589; U.S. Pat. No. 1,921,614; U.S. Pat. No. Des. 316,307; U.S. Pat. No. Des. 319,705; and U.S. Pat. No. 4,931,775.\nWhile these devices fulfill their respective, particular objectives and requirements the aforementioned patents do not disclose a new decorative ornamental assembly. The inventive device includes a pot assembly including pot members each having bottom and side walls and each also having an open top; and also includes a light-emitting assembly including candle-like light members being securely disposed in the pot members and extending through the tops of the pot members; and further includes poinsettia-like ornaments being disposed upon the pot members; and also includes support structures being attached to the pot members.\nIn these respects, the decorative ornamental assembly according to the present invention substantially departs from the conventional concepts and designs of the prior art and in so doing provides an apparatus primarily developed for the purpose of providing easy to use and colorful decorations about Christmas trees and on building structures.\nIn view of the foregoing disadvantages inherent in the known types of poinsettia-like ornaments now present in the prior art, the present invention provides a new decorative ornamental assembly construction wherein the same can be utilized for providing easy to use and colorful decorations about Christmas trees and on building structures.\nThe general purpose of the present invention, which will be described subsequently in greater detail, is to provide a new decorative ornamental assembly which has many of the advantages of the poinsettia-like ornaments mentioned heretofore and many novel features that result in a new decorative ornamental assembly which is not anticipated, rendered obvious, suggested, or even implied by any of the prior art poinsettia-like ornaments, either alone or in any combination thereof.\nTo attain this, the present invention generally comprises a pot assembly including pot members each having bottom and side walls and each also having an open top; and also includes a light-emitting assembly including candle-like light members being securely disposed in the pot members and extending through the tops of the pot members; and further includes poinsettia-like ornaments being disposed upon the pot members; and also includes support structures being attached to the pot members.\nThere has thus been outlined, rather broadly, the more important features of the invention in order that the detailed description thereof that follows may be better understood, and in order that the present contribution to the art may be better appreciated. There are additional features of the invention that will be described hereinafter and which will form the subject matter of the claims appended hereto.\nIn this respect, before explaining at least one embodiment of the invention in detail, it is to be understood that the invention is not limited in its application to the details of construction and to the arrangements of the components set forth in the following description or illustrated in the drawings. The invention is capable of other embodiments and of being practiced and carried out in various ways. Also, it is to be understood that the phraseology and terminology employed herein are for the purpose of description and should not be regarded as limiting.\nAs such, those skilled in the art will appreciate that the conception, upon which this disclosure is based, may readily be utilized as a basis for the designing of other structures, methods and systems for carrying out the several purposes of the present invention. It is important therefore, that the claims be regarded as including such equivalent constructions insofar as they do not depart from the spirit and scope of the present invention.\nFurther, the purpose of the foregoing abstract is to enable the U.S. Patent and Trademark Office and the public generally, and especially the scientists, engineers and practitioners in the art who are not familiar with patent or legal terms or phraseology, to determine quickly from a cursory inspection the nature and essence of the technical disclosure of the application. The abstract is neither intended to define the invention of the application, which is measured by the claims, nor is it intended to be limiting as to the scope of the invention in any way.\nIt is therefore an object of the present invention to provide a new decorative ornamental assembly which has many of the advantages of the poinsettia-like ornaments mentioned heretofore and many novel features that result in a new decorative ornamental assembly which is not anticipated, rendered obvious, suggested, or even implied by any of the prior art poinsettia-like ornaments, either alone or in any combination thereof.\nIt is another object of the present invention to provide a new decorative ornamental assembly which may be easily and efficiently manufactured and marketed.\nIt is a further object of the present invention to provide a new decorative ornamental assembly which is of a durable and reliable construction.\nAn even further object of the present invention is to provide a new decorative ornamental assembly which is susceptible of a low cost of manufacture with regard to both materials and labor, and which accordingly is then susceptible of low prices of sale to the consuming public, thereby making such decorative ornamental assembly economically available to the buying public.\nStill yet another object of the present invention is to provide a new decorative ornamental assembly which provides in the apparatuses and methods of the prior art some of the advantages thereof, while simultaneously overcoming some of the disadvantages normally associated therewith.\nStill another object of the present invention is to provide a new decorative ornamental assembly for providing easy to use and colorful decorations about Christmas trees and on building structures.\nYet another object of the present invention is to provide a new decorative ornamental assembly which includes a pot assembly including pot members each having bottom and side walls and each also having an open top; and also includes a light-emitting assembly including candle-like light members being securely disposed in the pot members and extending through the tops of the pot members; and further includes poinsettia-like ornaments being disposed upon the pot members; and also includes support structures being attached to the pot members.\nStill yet another object of the present invention is to provide a new decorative ornamental assembly that provides a novel ornamental appearance.\nEven still another object of the present invention is to provide a new decorative ornamental assembly that resemble life-like potted poinsettias.\nThese together with other objects of the invention, along with the various features of novelty which characterize the invention, are pointed out with particularity in the claims annexed to and forming a part of this disclosure. For a better understanding of the invention, its operating advantages and the specific objects attained by its uses, reference should be made to the accompanying drawings and descriptive matter in which there are illustrated preferred embodiments of the invention."}
{"text": "This invention relates to an image processing device for converting a scanned image into a monochrome image, and more particularly relates to an image processing device for processing a half tone between a black level and a white level.\nA conventional device converts a scanned image into a black and white image in order to process and to transmit the scanned image easily. Japanese laid open patent application Ser. No. 62-277857 published on Dec. 2, 1987 discloses a circuit for converting the scanned image into a black and white image. The circuit samples and holds a peak level of the scanned image signal in order to determine a threshold level. The circuit compares the image signal with the threshold level and generates a corresponding image signal having only black and white levels.\nHowever, the conventional device cannot process half tones between black and white levels, because the half tone levels fall into one of the black and the white levels. Additional circuits for determining double or triple threshold levels are required in order to process the half tone levels."}
{"text": "FIG. 1 illustrates a conventional block diagram of a television (TV) system. The TV system mainly includes a TV controller 100, dynamic random access memory (DRAM) 120, and a cathode ray tube (CRT) 140. The TV controller 100 is capable of receiving a plurality of types of video sources 110, including signals of National Television Standards Committee (NTSC) TV, Phase Alternation Line (PAL) TV, and video signals of super extended graphics array (SXGA)/extended graphics array (XGA)/video graphics array (VGA). The DRAM 120 is necessarily required to store frame data generated by the video sources 110 in FIG. 1. The TV controller 100 then properly displays frame data in the DRAM 120 to the CRT 140. The CRT 140 displays the video signal according to vertical synchronous (VSYNC) and horizontal synchronous (HSYNC) signals (not shown in FIG. 1) generated by the TV controller 100.\nHowever, the DRAM 120 increases the total manufacturing cost of the display system. Further, the CRT 140 cannot change VSYNC and HSYNC signals dramatically. The VSYNC and HSYNC signals can only vary their frequencies under vendor's specific specification, such as the specification of VSYNC signal within 5% and HSYNC signal within 2%. Otherwise the CRT 140 will be seriously damaged or the CRT video display is seriously distorted on the CRT 140. The VSYNC signal is 60 Hz and 50 Hz for NTSC and PAL standards, respectively. SXVGA/XGA/VGA video signals may support a plurality of kinds of VSYNC operation frequencies, preferably ranging from 60 Hz to 85 Hz.\nIn the prior art, the TV controller 100 separates the input frame data from the output frame data by buffering frame data into the DRAM 120, so the output video signals composed of the output frame data can merely be controlled under approximate fixed frequencies of VSYNC and HSYNC signals to stably display the video signals. Specifically, when a different video source is selected into the TV controller 100 or when the display mode is changed, the TV controller 100 may distort video signals on the CRT 140 or even damage the CRT 140 due to the VSYNC and HSYNC signal variations.\nFIG. 2 shows a block diagram of the conventional TV controller in FIG. 1. The TV controller 100 comprises an output phase-locked loop (PLL) 200 and a horizontal and vertical (HV) generator 220. The output PLL 200 receives a fixed input clock 210 and then outputs Output_CLK signal. The HV generator 220 generates VSYNC and HSYNC signals according to Output_CLK signal. The conventional display controller processes image frames based on pixels. Persons skilled in the art considers CRT TV cannot be implemented without the DRAM 120 since the output VSYNC signal must change significantly to cause the CRT to be seriously damaged or cause the CRT video display to be seriously distorted when the input TV video source is changed. It's unacceptable by the manufactures of the display systems, and especially damaging CRT TV is too dangerous.\nAs aforementioned, either CRT or LCD controller, of which conventional PLL receives a fixed input clock, cannot afford the demand of different video sources and DRAM increases the manufacturing cost of display systems. Consequently, there is a need to develop a display controller to provide the display system with rapid adjustment of VSYNC and HSYNC signals for reducing the cost and improving stability thereof."}
{"text": "Human-computer interactions have progressed to the point where humans can control computing devices, and provide input to those devices, by speaking. Computing devices employ techniques to identify the words spoken by a human user based on the various qualities of a received audio input. Such techniques are called speech recognition or automatic speech recognition (ASR). Speech recognition combined with language processing techniques may allow a user to control a computing device to perform tasks based on the user's spoken commands. Speech recognition may also convert a user's speech into text data, which may then be provided to various textual based programs and applications.\nComputers, hand-held devices, telephone computer systems, kiosks, and a wide variety of other devices to improve human-computer interactions may use speech recognition."}
{"text": "The present invention is directed to bicycles and, more particularly, to a switch designation apparatus for adjusting a position of a bicycle control unit.\nBicycle transmissions that comprise front and rear derailleurs are well known. The front derailleur shifts a chain among a plurality of front sprockets that are coaxially mounted to the pedal crank shaft, and the rear derailleur shifts the chain among a plurality of rear sprockets that are coaxially mounted to the rear wheel. It is also known to use electric motors to operate the front and rear derailleurs. Such a system is shown in Japanese Kokai Patent Application No. 2001-267002.\nBicycle transmissions formed inside the bicycle wheel hub (often called internal hub transmission) also are well known. Such transmissions include a planetary gear mechanism mounted inside the hub for communicating rotational force from a drive sprocket to the wheel hub shell through a plurality of transmission paths representing different gear ratios. It is also known to use electric motors to select the desired transmission paths. Such a system is shown in Japanese Patent Application No. 3-231006.\nAn apparatus for controlling the operation of a bicycle transmission having electrically operated derailleurs typically comprises one or more shift control devices mounted to one or more sides of the bicycle handlebar. For example, a shift control device for controlling the operation of the front derailleur may be mounted to the left side of the handlebar, and a shift control device for controlling the operation of the rear derailleur may be mounted to the right side of the handlebar. Each shift control device may comprise a switch for upshifting its corresponding derailleur and a switch for downshifting its corresponding derailleur, wherein the upshift and downshift switches are arranged vertically in a row. A shift control device for an internal hub transmission may comprise a switch for upshifting the internal hub transmission and a switch for downshifting the internal hub transmission, wherein both switches are mounted horizontally in a row on the right side of the handlebar.\nWhether the switches are arranged vertically or horizontally, the upshift and downshift operations of the switches are fixed. However, some riders may prefer a reversal of the physical arrangement of the upshift and downshift switches. In other words, while one rider may prefer the upshift switch located above the downshift switch, other riders may prefer the downshift switch to be located above the upshift switch. With conventional shift control devices, this would require mounting the shift control device upside down or rewiring the device. Mounting the device upside down would result in any printing on the device also being upside down, and the switches may be oriented in a non-ergonomic manner. Rewiring the device would be very complex and troublesome."}
{"text": "As known from the prior art, in a dryer section of a paper machine, single-wire draw or twin-wire draw or various combinations of the same are employed. In single-wire draw, so-called \"normal groups with single-wire draw\" are used, in which drying cylinders are arranged in an upper row and reversing suction cylinders or rolls are arranged in a lower row below the upper row of drying cylinders. In this case, the dryer section is open toward the bottom, for example, to enable easy removal of broke. In so-called \"inverted groups with single-wire draw\", the drying cylinders are arranged in a lower row and the reversing suction cylinders or rolls are arranged in an upper row above the lower row of reversing suction cylinders or rolls. In this case, it is a recognized drawback that the inverted groups are closed toward the bottom and, for example, there are problems associated with the removal of broke because it cannot be arranged to take place by the force of gravity (which is how it is done in \"normal groups with single-wire draw\").\nIn the past, in the dryer section of a paper machine, so-called breaker stacks (intermediate calenders) were employed quite commonly, which stacks are usually arranged between groups with twin-wire draw (which will be discussed below). These breaker stacks have been largely abandoned, partly because of runnability problems, and partly because their us causes an increase in the length of the dryer section in the machine direction since, between a breaker stack and the wire groups preceding and following it, it was usually necessary to provide relatively long unsupported draws of the web. Such draws are susceptible to fluttering and web breaks and are also problematic in view of the threading of the web.\nIn recent years, some breaker stacks arranged inside dryer sections have been suggested, which stacks are formed between heated drying cylinders and particular calender rolls. With respect to this prior art, reference is made, as an example, to the U.S. Pat. No. 5,127,168, to published German Patent Application No. DE 4,407,405 A1 (which corresponds to English-language Canadian Patent Application No. 2,143,912), and to a press section marketed by Black Clawson--Kennedy Inc. under the trademark \"HYDRA NIP\".TM.. It is not known if a patent or patent application in respect of the latter press section has been filed or issued. With respect to the prior art, reference is made further to Japanese Patent Application No. 56040/1992 (published application No. 222691/1993), in the name of Mitsubishi Heavy Industries Ltd.\nIn FIG. 5 in U.S. Pat. No. 5,127,168, an arrangement is illustrated in which a first nip of intermediate calendering is formed in connection with the first cylinder in the first group with twin-wire draw in a multi-cylinder dryer, and a water-receiving felt of the press section of the paper machine is passed through this first nip of intermediate calendering. A second nip of intermediate calendering is formed in connection with the first lower cylinder in the second group with twin-wire draw. It is a drawback of the arrangement of intermediate calendering known from this U.S. patent that in the group gap in which an intermediate calendering nip is used, the web will have a very long unsupported draw, and further, it is a disadvantage that the group gap makes the dryer section remarkably longer in the machine direction. The U.S. patent also does not suggest any solutions for carrying out intermediate calendering in modern dryer sections in which groups with single-wire draw and closed draw of the web are applied.\nIn FIG. 1 of the German Patent Application No. DE 4,407,405, an arrangement of intermediate calendering is illustrated in which an intermediate calendering nip is arranged in connection with the last drying cylinder in an inverted group with single-wire draw. The arrangement of intermediate calendering in accordance with the German patent application involves the drawback that the group gaps in which intermediate calendering is applied become relatively long and spacious in the machine direction and thereby unduly increase the length of the dryer section in the machine direction. Moreover, in the German patent application, an arrangement of intermediate calendering is described as being applied in connection with an inverted group with single-wire draw and in group gaps between such an inverted group and a normal group with single-wire draw only. As such, the German patent application does not suggest any solution for intermediate calendering in dryer sections in which exclusively groups with single-wire draw or various combinations of such groups and groups with twin-wire draw or combinations of such groups and so-called hybrid groups are employed.\nIn particular in connection with modernizations of paper machines, for example in order to increase their running speed, it is also necessary to increase the drying capacity of the dryer section. The prior art arrangements of intermediate calendering mentioned above, with the exception of that described in U.S. Pat. No. 5,127,168, are not suitable for modernizations of paper machines, because they increase the length of the dryer section and/or reduce the drying capacity."}
{"text": "1. Field of the Invention\nThe invention relates to thermally conductive pads and, particularly, to a thermally conductive pad with an array of carbon nanotubes and a method for making the same.\n2. Description of Related Art\nSince the report of their discovery in 1991 by Iijima, carbon nanotubes (CNTs) have been extensively studied for their structural, physio-chemical, mechanical, electrical, and electromechanical properties. Many potential technological applications have been proposed including hydrogen storage, nanoelectronic devices, field emission displays (FED), field emission microscopy (FEM), chemical sensors, and so on. Carbon nanotubes are desirable, at least in part, due to their unique electrical and mechanical properties. Currently, the study of carbon nanotube/polymer composite materials is gaining significant attention. Such nanotube-reinforced composite materials have broad applications because of their good antistatic performance, microwave absorbing capability, electromagnetic shielding ability, and so on.\nGenerally speaking, composite materials with carbon nanotubes have been prepared via an in-situ polymerization process, a solution compounding process, and/or a melt compounding process.\nCarbon nanotubes arrays have excellent mechanical strength and good heat conductivity and/or heat diffusivity, as such they have been used to produce heat conductive materials and reinforced composite materials. However, the density of the carbon nanotubes in the composite material effects the mechanical strength and the heat conductivity of the composite material.\nCurrently, the technology for chemical vapor deposition (CVD) method for producing/making carbon nanotube arrays has reached a mature state. However, the carbon nanotube arrays directly produced or grown by the CVD method have a density of less than 0.01 grams per cubic centimeter (g/cm3). Thus, the carbon nanotube arrays are not tightly bounded to each other, and the spacings between the carbon nanotubes are in the order of several times the diameter of the carbon nanotubes. Furthermore, the density of the carbon nanotube arrays directly grown by the CVD method is difficult to change and/or control. As such, composite materials containing low-density carbon nanotube arrays have poor thermal conductivity. Thereby, the composite materials may not achieve the desirable effects required of thermally conductive composite materials.\nFurther, a thermally conductive pad can be obtained by cutting the above-described thermal conductive composite material. However, the thermally conductive pad has a low coefficient of thermal conductivity. Thus, the thermally conductive pad is restricted in its use in the thermal conductive field.\nWhat is needed, therefore, is to provide a thermally conductive pad with a high-density array of carbon nanotubes, and a method for producing the same that is very simple and efficient and the density of carbon nanotubes therein is easily controlled during the production process."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for preparing (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo [1,2-a][1,2,4]triazolium-6-yl)]thio-6-[(R)-1-hydroxyethyl]-1-methyl-carbap enem-3-carboxylate represented by the following formula (I) ##STR3## or a pharmacologically acceptable salt thereof.\n2. Description of the Prior Art\nSince the discovery of an antibiotic substance \"thienamycin\" from the nature [U.S. Pat. No. 3,950,357; J. Am. Chem. Soc., 100, 313 (1987)], many researches have been carried out to develop carbapenem antibiotics. As a result of the extensive study, there has been developed imipenem (INN) which is a practically available antibacterial agent for clinical use.\nRecently, Kumagai et al. proposed (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo [1,2-a][1,2,4]triazolium-6-yl)]thio-6-[(R)-1-hydroxyethyl]-1-methyl-carbap enem-3-carboxylate of formula (I) above as carbapenem antibiotic. See U.S. Pat. Nos. 4,866,171; 4,925,836; and 4,925,935 (the entire specifications of these patents are incorporated herein by reference). This compound is carba-2-penem-3-carboxylic acid compound in which a .beta.-methyl group is introduced at 1-position and 6,7-dihydro-5H-pyrazolo[1,2-a][1,2,4]triazolium-6-yl thio group which is represented by the following ##STR4## is introduced at 2-position of the carbapenem skelton. Owing to the presence of these specific substitutions, the compound of formula (I) has superior antibacterial activity with excellent chemical and physicochemical stabilities in living bodies, and is extremely stable against dehydropeptidase (DHP) known as a kidney enzyme.\nThe only synthesis for preparation of (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo [1,2-a][1,2,4]triazolium-6-yl)]thio -6-[(R)-1-hydroxyethyl]-1-methylcarbapenem-3-carboxylate of formula (I) heretofore known in Kumagai et al. is from an imino intermediate. That is, according to Kumagai et al., the comound of formula (I) can be prepared by reacting a compound represented by formula (II): ##STR5## wherein R.sup.1 is a carboxyl protecting group, and R.sup.a is an acyl group,\nwith a mercapto reagent represented by formula (VII): ##STR6## wherein R.sup.b is an amino protecting group, to give a compound represented by formula (VIII): ##STR7## wherein R.sup.1 and R.sup.b have the same meanings as above, and subjecting the compound of the formula (VIII) to removal of the protecting groups R.sup.1 and R.sup.b to give the carbapenem compound of the formula (IX) ##STR8## and then, reacting the resulting compound of formula (IX) with formimidic acid ester to give the carbapenem compound of formula (I).\nOne presumed course of reaction in the formation of the compound of formula (I) by the reaction of the compound of formula (IX) with the formimidic acid ester derivative is that by the reaction of the compound (IX) with the formimidic acid ester derivative, (1R,5S,6S)-6-[(R)-1-hydroxyethyl]-1-methyl-2-[(1,2-di-iminomethyl)-4-pyraz olidinyl]thiocarbapenem-3-carboxylic acid of the following formula ##STR9## is formed as an intermediate, and this imino intermediate is transitory compound and undergoes cyclization reaction to form the compound of formula (I).\nFinding a method for preparing the compound (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo [1,2-a][1,2,4]triazolium-6-yl)]thio-6-[(R)-1-hydroxyethyl]-1-methylcarbape nem-3-carboxylate of formula (I) through a method other than from the imino intermediate has been an objective of the research of the present inventors and their assignee since the discovery of the useful therapeutic indications of the compound (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo [1,2-a][1,2,4]triazolium-6-yl)]thio-6-[(R)-1-hydroxyethyl]-1-methyl-carbap enem-3-carboxylate. Despite the efforts of the researches of the assignee to find a way to make (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo [1,2-a][1,2,4]triazolium-6-yl)]thio-6-[(R)-1-hydroxyethylethyl]-1-methyl-c arbapenem-3-carboxylate without going through the transitory intermediate, efforts heretofore have been unsuccessful. One apparent method for making (1R,5S,6S)-2-[(6,7-dihydro-5H-pyrazolo[1,2 -a][1,2,4]triazolium-6-yl)]thio-6-[(R)-1-hydroxyethyl]-1-methylcarbapenem- 3-carboxylate is disclosed in Merck European Patent Application 168,707 (1986), that is discussed in detail in Kumagai et al. U.S. Pat. No. 4,866,171 at column 1, lines 22-49. There is a picture representation of the compound 6,7-dihydro-6-mercapto-5H-pyrazolo[1,2-a]-[1,2,4-triazolium halide which is a key compound for the synthesis in the European application. Nowhere in the European application is there any disclosure of the synthesis of a 6,7-dihydro-6-mercapto-5H-pyrazolo[1,2-a][1,2,4]triazolium halide nor is any method for the synthesis of a 6,7-dihydro-6-mercapto-5H-pyrazolo[1,2-a]-[1,2,4]triazolium halide available within the skill of the art."}
{"text": "This application claims the benefit under 35 U.S.C. 119 of Great Britain application number GB 0012229.1, filed May 19, 2000.\nThe invention relates to the use of microscopic nematodes such as C. elegans in functional high throughput in vivo assays suitable for the detection of inhibitors or activators of intestinal lipid uptake. Compounds identified as modulators of intestinal lipid uptake using the assays of the invention may provide lead compounds for the development of pharmaceutical agents useful in the treatment of diseases of the human or animal (e.g. mammal) body, and in particular diseases and disorders of human and/or animal metabolism, fat handling and/or fat storage, such as obesity, impaired fat metabolism and other related diseases such as diabetes type II and cardiovascular diseases.\nWhen fat reaches the intestines in vertebrates, the pancreatic lipase hydrolyses the triglycerides into smaller components designated free fatty acids and monoglycerides (mainly 2-monoacylglycerols). Fatty acids are long-chain hydrocarbon molecules containing a carboxylic moiety at one end. The numbering of carbons in fatty acids begins with the carbon of the carboxylate group.\nMetabolically, fatty acids are important energy substrates because of their high calorific content. In a typical diet of Western developed countries, approximately 30-40% of the dietary calories are derived from lipids, mainly in the form of di- and triglycerides. The linkage between excessive dietary lipid consumption and several common pathophysiologic disorders, including heart disease, obesity and diabetes and cancer, has been widely documented (Watts, et al., Am J Clin Nutr 64, 202-9 (1996); Storlien et al., Science 237, 885-8 (1987)).\nThere are three major roles in the body for the free fatty acids:\n1) as the components of more complex membrane lipids.\n2) as the major components of stored fat in the form of triacylglycerols.\n3) Metabolism of fatty acids by xcex2-oxidation is also the major source of ATP as energy for most organisms, especially for mammalian cardiac muscle.\nUntil recently it was considered that the adsorption of fatty acids into the body during digestion was through passive diffusion rather than the active transport process, as was known for carbohydrates and amino acids. Presently at least five plasma membrane proteins have been identified and proposed as candidates for fatty acid transporters thus far. These include, but are not restricted to:\nPlasma Membrane Fatty Acid Binding Protein (FABPpm),\nFatty Acid Translocase (FAT)\nCaveolin, a 22-kDa fatty acid binding protein\nRenal 56-kb FABP\nFatty Acid Transport Protein (FATP)\nAn overview of these membrane proteins has been published by Yuen Hui and David A. Bernlohr, Bioscience 2, 222-231 (1997).\nThe expression of FATP is regulated by certain transcription factors, such as the xe2x80x9cPPARxe2x80x9d (peroxysome proliferator activated receptor)-transcription factors the xe2x80x9cRXRxe2x80x9d (Retinoid X receptor)-transcription factors, and similar factors. Therefore, activators of these receptors, respectively fibrates or antidiabetic thiazolidinedione and retinoic acid, can increase FATP expression. One of the six known human FATPs, FATP4, has recently been shown to possess the functional characteristics (presence and absence is correlated with increase or decrease in fatty acid uptake) and cellular location (highly expressed in the microvilli of intestinal enterocytes) that would be required in a major intestinal fatty acid transport protein (Stahl A, et al. Mol Cell Biol 4,299-308 (1999)). It is highly probable that the expression and activity of the other transporter proteins is regulated too.\nNucleotide sequences encoding for and protein sequences of these fatty acid transporter proteins, both the human proteins and their C. elegans homologues, can be found in publicly accessible sequence databases, such as GenBank (accessible at the National Center for Biotechnology Information website, http://www.ncbi.nlm.nih.gov/PubMed/) and the C. elegans database of the Sanger Centre, UK (accessible at the Sanger Centre website, http://wormbase.sanger.ac.uk/). Some examples of sequences and designation numbers are:\nAlthough less well documented, proteins involved in the uptake of other lipids have been described in literature, and sequences of these proteins and genes have been published. For example, it has been shown that ABC transporters and more particularly the ABC transporter ABCB1 (ABC8) plays an important role in the regulation of Cholesterol and phospholipid transport (Klucken et al., Proc. Natl. Acad. Sci. USA, 2000, 97:817-822). Delivery of lipids, and more particularly sterols and cholesterol is done by the scavenger receptor-BI (Stangl et al., J. Biol. Chem., 1999, 274:32692-32698). Izzat et al. describe other putative targets to reduce the intestinal cholesterol uptake in the Journal of Pharmacology and Experimental Therapeutics 2000, 293:315-320. Such targets include bile acid transporters and HMG-CoA reductase. Cholesterol is taken up by the gut membrane without the involvement of a transporter. Compounds that interact stoichiometrically with the cholesterol in the intestinal lumen would also reduce cholesterol uptake. Other transporter proteins can be found in literature.\nControlling the uptake of lipids in the intestines would allow the treatment of obesity as well as the treatment of some related diseases such as diabetes mellitus, cardiovascular diseases such as arteriosclerosis, hypertension, stroke and certain forms of cancer. Enhancing lipid uptake in other tissues may also have specific therapeutic applications. For example, enhanced uptake of fatty acids by the skin will result in improved cosmetics.\nThe present inventors describe herein a functional assay to measure the uptake of lipids in vivo in a real intestine environment using small microscopic nematodes, such as Caenorhabditis elegans, as an animal model. Specific applications of this method are also described.\nTherefore, in accordance with a first aspect of the invention there is provided a method of assaying lipid uptake in microscopic nematode worms which comprises:\nincubating the said microscopic nematode worms in the presence of a probe molecule comprising a lipid moiety linked to a signal generating label; and\ndetermining the amount of probe molecule taken up by the said microscopic nematode worms by detecting a signal generated by the label part of the probe molecule.\nIn the assays of the invention, the nematode worms are incubated to a medium containing the probe molecule. Upon such incubation, the probe is taken up into the gastrointestinal tract of the nematode (e.g. by pharynx pumping), and in particular into the gut (lumen) of the nematode.\nFrom the gut lumen, the probe may then pass from the gut lumen through the wall of the gastrointestinal tract (i.e. of the gut) into the body of the nematode, where it may concentrate in specific cells or tissues, such as the gut granules and/or other cells and tissues that store or handle lipids.\nAccording to the invention, this passing of the probe molecule(s) from the gut lumen of the nematode into the body of the nematode is used as an in vivo model for lipid transport across biological membranes or barriers, not just in nematodes, but also in higher multicellular organisms, such as vertebrates, mammals and even humans. (In this respect, it should also be noted that the passing of the probe molecule(s) from the gut lumen of the nematode into the body of the nematode may not just be used as a model for lipid transport across the wall of the gastrointestinal tract, but generally as a model for lipid transport across any biological membrane and/or barrier, including epithelial cells/cell layersxe2x80x94such as the walls of blood vesselsxe2x80x94and cell membranes.)\nThe invention is not particularly limited as to the mechanism or pathway via which said lipid transport may take place (i.e. both in nematode used in the assay of the invention as well as in the higher organism for which the nematode serves as an in vivo model) and may for instance include mechanisms such as active transport (which is preferred) and (passive) diffusion, or a combination thereof. In one embodiment, the invention may even be used to determine if, under the conditions of the assay (i.e. as used as a model), a probe or lipid is transported by an active transport mechanism or by a passive transport mechanism (e.g. by diffusion).\nThis may also depend on the specific type of lipid or probe involved (i.e. on the type of lipid moiety present in the probe), andxe2x80x94for the assay of the inventionxe2x80x94on the specific worm strain used, on genetic factors influencing the worm strain used (such as gene suppression, e.g. as induced by RNAi methodology) and/or on the induction or suppression of specific pathways in the nematode(s) used, e.g. as a result of exposure of the worms to compounds that induce or suppress (the expression of) certain pathways involved in lipid transport and/or of specific genes and/or enzymes that are involved in such pathways (as further described below.)\nAs will be appreciated, to measure (active or passive) transport of a probe molecule from the gut lumen into the body of the nematode in vivo, it would in principle be necessary to measure, over a relevant period of time, either the amount and/or concentration of probe in the gut lumen (and/or any changes therein), the amount and/or concentration of probe in the body of the nematode (and/or any changes therein), or both. However, this would be very elaborate, would require examination of individual wormsxe2x80x94e.g. using microscopy techniquesxe2x80x94and even so would probably still not afford quantitative and/or statistically relevant results. Also, any such measurement would not be suitable for automation and/or for screening at medium to high throughput.\nThus, it is one of the general objects of the present invention to provide an assay for measuring lipid uptake/transport in vivo, which assay can serve as an in vivo model for lipid uptake/transport across biological membranes or barriers in (higher) multicellular animals, and which assay can be susceptible toxe2x80x94i.e. can be configured forxe2x80x94automation and/or to screening at medium to high throughput.\nMore specifically, it is an object of the invention to provide such an assay using nematode worms as a model organism.\nOne particular purpose of the invention is to provide such an assay which can be used for determining the influence on such lipid uptake/transport of small molecules, of induction or suppression of pathways involved in lipid uptake/transport (e.g. by prolonged exposure to inducing or suppressing factors), and/or of genetic factors such as mutations or RNAi-induced gene suppression.\nIn this way, the assay of the invention could be used to screen libraries of chemical compounds for small molecules and/or for (other) factors that can influence lipid uptake/transport across biological membranes/barrier, which small molecules or factors could then be used (as a starting point) in the development of compounds for use in compositions for increasing or decreasing lipid uptake/transport across biological membranes/barriers in multicellular organisms, including vertebrates such as mammals and humans, and more generally for influencing or altering fat and/or lipid handling or storage by such multicellular organisms, which can have relevance for a number of disease areas in animals and humans, including but not limited to obesity, diabetes and cardiovascular diseases (and in particular those relating to cholesterol handing and metabolism).\nIn addition to screening chemical libraries, the assay of the invention could also be used in drug development, e.g. in hits-to-leads chemistry or lead development, but also in genetic screens, gene discovery techniques, target validation techniques and/or other (functional) genomics techniques.\nIn a specific embodiment, the assays of the invention may also be used to measure/determine the uptake of specific components of food or food compositions (e.g. for human or animal use), such as fatty acids, fats, oils cholesterol-like compounds and/or other lipids (e.g. as defined below), as well as of compounds that are intended to replace fats, oils and/or lipids (for instance in food compositions for human consumption), and/or for additivesxe2x80x94e.g. for food compositionsxe2x80x94that are intended to lower cholesterol in the animal or person that consumes said composition. In this specific aspect of the invention, the assay may be used as a model for the uptake of such compounds from the gastrointestinal tract of a higher animal, such as an animal or human, e.g. to determine components which show increased or reduced uptake from the g.i. tract (i.e. compared to a reference). Thus, the assays of the invention may also be used in the development of components/additives for food, such as dietary compositions for human consumption, (specialized) nutritional compositions, and/or infant formula, or components for use therein.\nThe invention achieves the objects referred to above by providing the assay techniques described herein."}
{"text": "Distributed asynchronous message driven systems are inherently complex when it comes to handling partial failures. Data flow diagrams can be used to describe and design the staged computation in distributed message driven systems by showing how messages flow between components. Data flow diagrams can drive the automatic synthesis of application executables, i.e. source code can be generated from data flow diagrams.\nData flow diagrams, however, have limitations when it comes to partial failure handling. In a data flow, a failure manifests as a fault message. The functional path of a data flow can be instrumented with junctions for fault messages. However, this approach is unfeasible for handling failures because multiple different failures can occur at any stage (partial failure) and can affect parallel data flows. Furthermore, multiple instances of a data flow can be active at the same time. For this reason, capturing which part in which instance of a data flow is affected by which kind of failure would require diagramming all combinations of data paths under faults. Such an approach is simply unfeasible in most situations.\nDistributed asynchronous message driven systems are also inherently complex to debug. Traditional stack-based debuggers are not typically useful in distributed asynchronous environments since local stacks are meaningless. Similarly, setting traditional breakpoints on source code lines is also of little value in such an environment since multiple instances of the same lines of code are running concurrently with each other."}
{"text": "This invention concerns semiconductor processing methods for filling structural gaps and integrated circuitry having filled structural gaps.\nIn semiconductor wafer processing, device structures are formed on a semiconductor substrate. FIG. 6 illustrates a wafer fragment 600 comprising a bulk substrate 45 having a pair of wordline constructions 60 formed thereon. A gate oxide 15, gate conductor 20, cap 25, and spacers 30 comprise the individual structural components of wordline constructions 60. Often, it is desirable to form a layer of insulation material over the device structures formed on a semiconductor wafer, such as substrate 45. An insulation layer 640 on wafer portion 600 may be formed by depositing borophosphosilicate glass (BPSG) and reflowing the deposited insulation material to densify and planarize insulation layer 640.\nBefore forming insulation layer 640, a gap exists in areas where substrate 45 is not covered by device structures, like wordline constructions 60. Such a gap may expose substrate 45 to boron and/or phosphorous diffusion from the BPSG in insulation layer 640. A barrier layer 635 is often formed over device structures prior to formation of insulation layer 640 to counteract such boron and/or phosphorous diffusion. Unfortunately, barrier layer 635 forms a reduced gap within the original gap that has a higher aspect ratio than the original gap. When the aspect ratio of a gap between device structures is sufficiently high, it may cause formation of voids, such as a void 650, in insulation layer 640 during deposition. It is highly desirable to form insulation layer 640 such that it fills the reduced gap between wordline constructions 60 without formation of void 650.\nAccordingly, barrier layer 635 addresses the problem of boron and/or phosphorous diffusion, but exacerbates the problem of void formation. Thus, a need exists to provide a method for protecting against boron and/or phosphorous diffusion while reducing the likelihood of forming voids in structural gaps. Otherwise, semiconductive substrates may either suffer defects resulting from boron and/or phosphorous diffusion or defects resulting from formation of voids.\nWhile motivated from this perspective, the artisan will appreciate other applicabilities, with the invention only being limited by the accompanying claims appropriately interpreted in accordance with The Doctrine of Equivalents.\nAccording to one aspect of the invention, a semiconductor processing method for filling structural gaps includes depositing a substantially boron free silicon oxide comprising material at a first average deposition rate over an exposed semiconductive material in a gap between wordline constructions and at a second average deposition rate less than the first average deposition rate over the wordline constructions. A reduced gap having a second aspect ratio less than or equal to a first aspect ratio of the original gap may be provided. A layer of a boron containing silicon oxide material may be deposited on the substantially boron free silicon oxide material and over the wordline constructions. The boron containing silicon oxide material may be deposited in situ in the same chamber.\nAccording to another aspect of the invention, an integrated circuit includes a pair of wordline constructions separated by a gap therebetween in areas where the wordline constructions do not cover an underlying semiconductive substrate. A layer of substantially boron free silicon oxide material has a first thickness over the substrate within the gap and has a second thickness less than the first thickness over the wordline constructions. A layer of boron containing silicon oxide material is over the substantially boron free silicon oxide comprising material.\nOther aspects are, of course, contemplated."}
{"text": "Washing machine appliances generally include a cabinet which receives a tub for containing wash and rinse water. A wash basket is rotatably mounted within the wash tub. A drive assembly is coupled to the wash tub and configured to rotate the wash basket within the wash tub in order to cleanse articles within the wash basket. Upon completion of a wash cycle, a pump assembly can be used to rinse and drain soiled water to a draining system.\nWashing machine appliances include vertical axis washing machine appliances and horizontal axis washing machine appliances, where “vertical axis” and “horizontal axis” refer to the axis of rotation of the wash basket within the wash tub. Irrespective of the axis, washing machine appliances may include multiple corners or support feet on which a particular appliance rests. In certain situations, a washing machine appliance may become tilted from its support feet and fall over. For example, an extreme imbalance (e.g., caused by an inappropriately sized or distributed load within the washing machine appliance) may rock the cabinet until the washing machine falls on one of its sides. As another example, the support surface on which the cabinet rests (e.g., the floor) may crack or break, causing the washing machine to become imbalanced and fall. Once the washing machine appliance has fallen, one or more moving or electrified portions of the washing machine appliance may become exposed to the user. This may include the drive assembly, which rotates the basket within the tub. Aside from potentially damaging the washing machine appliance itself, exposure of such portions may risk damaging or injuring nearby appliances or users.\nSome existing washing machine appliances may include features for confirming that the cabinet of a washing machine appliance is level before operations begin (e.g., during installation). However, typical existing appliances fail to provide any features for ensuring appliance or user safety once the washing machine appliance has fallen. In many cases, a user is required to turn off or unplug the washing machine appliance once the cabinet has moved onto its side. Furthermore, typical existing appliances fail to provide any features for preventing such unsafe conditions from occurring.\nAccordingly, a need exists for a washing machine appliance that includes features for addressing one or more of the above-identified issues. In particular, a washing machine appliance that includes features for minimizing the potential or risk of the appliance in a fallen state be advantageous."}
{"text": "1. Field of the Invention\nThis invention relates to control systems and, more particularly, to methods of driving dynamic and steady-state behavior of a process toward more optimum operating conditions.\n2. Discussion of Related Art\nMultivariable predictive control algorithms, such as DMCplus™ from AspenTech® of Aspen Technologies, Inc. or RMPCT from Honeywell International Inc., are a combination of calculations to drive the dynamic and steady-state behavior of a process toward a more optimum operating condition. The steady-state algorithm used in the control scheme is most commonly a linear program (LP), but sometimes is a quadratic program (QP). For small problems, understanding the LP or QP solution is relatively simple. Two-dimensional problems can be visualized on a paper and demonstrated to an operator to gain understanding of the process. With some detailed modeling background, engineers and well trained operators can understand medium-sized problems (less than 10 dimensions). However, larger, more interactive problems, often require offline simulation. This can take a significant amount of time to understand, even qualitatively.\nTypically, an operator of a multivariable predictive controller (MPC) can observe current constraints and may have access to an open loop model of the process. However, to fully understand the constraint set relief, the operator would need a detailed understanding of the process model and the ability to trace independent and dependent relationships through the model. For that, an offline simulation or analysis tool is required. Otherwise, the operator cannot know how much to change a constraint or which constraint is the next to become active.\nOne concept for an offline simulation uses a matrix pivot in which unconstrained manipulated variables (MVs) are swapped with constrained controlled variables (CV). The constraints become “independents,” and the unconstrained variables become “dependents.” The matrix pivot can be symbolized as follows:\n      [                                        y            1                                                            y            2                                ]    =                              [                                                    A                                            B                                                                    C                                            D                                              ]                ⁡                  [                                                                      x                  1                                                                                                      x                  2                                                              ]                    ⁢                          [                                                  x              1                                                                          y              2                                          ]        =                  [                                                            A                                  -                  1                                                                                                      -                                      A                                          -                      1                                                                      ⁢                B                                                                                        CA                                  -                  1                                                                                    D                -                                                      CA                                          -                      1                                                        ⁢                  B                                                                    ]            ⁡              [                                                            y                1                                                                                        x                2                                                    ]            \nHowever, this approach does not provide quantitative answers as to how much any operator change will affect the controller solution.\nThere is a need for a simple utility that can analyze past or current dynamic matrix control (DMC) solutions of any-sized problem, in real-time, to provide an operator meaningful, quantitative instructions for DMC controller constraint relief."}
{"text": "1. Field of the Invention\nThe present invention relates to a coating composition excellent in weatherability and abrasion resistance, and a resin molded article coated with the composition.\n2. Prior Art of the Invention\nIn general, polycarbonate resins are widely used as a transparent resin material due to their excellent impact resistance and transparency. However, polycarbonate resins have defects in that they are poor in abrasion resistance and solvent resistance and that the surface of a molded article produced therefrom is easily scratched and easily corroded with an organic solvent. In order to overcome these defects, there have been proposed a variety of methods in which such a molded article is coated with a thermosetting resin.\nIn order to improve the properties, abrasion resistance, etc., of transparent resin articles such as molded articles produced from the above polycarbonate resins, well known are resin articles which are coated, as an undercoat layer, with a coating composition an acrylic polymer and further coated, as an overcoat layer, with a coating composition comprising either a silicon compound or a partial hydrolysis product of a silicon compound and a colloidal metal oxide composed mainly of colloidal silica (TOKKAI 59-109528(1984), etc.). And, the undercoat layer of an acrylic polymer often contains an ultraviolet light absorber in order to prevent crack and a decrease in adhesion to another coating layer, which occur due to photo-oxidative deterioration caused on an article surface by ultraviolet light.\nHowever, such resin articles sometimes turn yellow, suffer cracks or have a coating layer peeled off under the influence of ultraviolet light, water, heat, etc., when used for a long period of time. These phenomena are caused for the following reason. When such resin articles are used for a long period of time, the ultraviolet light absorber migrates onto the coating layer surface and dissipates, i.e. it would not remain on the resin article surface, and the ultraviolet light absorption effect is gradually decreased.\nFurther, when there is employed a method such as an immersion method, a flow-coat method, or the like, in which an overcoating liquid for coating an undercoated resin article with an overcoat layer is repeatedly used, the ultraviolet light absorber container in the undercoat layer is eluted into the overcoating liquid when brought into contact with the overcoating liquid, and the concentration of the ultraviolet light absorber in the overcoating liquid gradually increases. As a result, the resultant overcoating liquid shows defectiveness in hardness, a decrease in adhesion properties, etc., and the pot life of such an overcoating liquid decreases."}
{"text": "Metallization is a key element in semiconductor technology. Metallization is in charge of current transport in and out of a semiconductor chip as well as removal of heat generated during operation of the semiconductor chip. Metal adhesion and barrier structures aim at providing adhesion between metallization and a support structure such as a semiconductor body and preventing metal atoms from diffusing from the metal structure into a semiconductor substrate. It is desirable to improve reliability of barrier and adhesion characteristics over a desired time span, to reduce damage of barrier characteristics due to defects and particles as well as to improve ability to screen barrier defects."}
{"text": "Compliant foil hydrodynamic fluid film radial bearings are currently being utilized in a variety of high speed rotor applications. These bearings are comprised of an inner rotating element, non-rotating compliant fluid foil elements that surround the rotating element, non-rotating compliant spring foil elements that surround the fluid foil element and a non-rotating cartridge element that surrounds and provides attachments for the foil elements. The annular space between the rotating element and the cartridge element is filled with fluid (usually air) which envelops the foils. The space between the smooth rotating element and the fluid foil is subdivided into a plurality of fluid-dynamic wedge channels. These wedge channels are formed by overlapping the fluid and spring foils to form ramps in the inner surfaces of the fluid foils. The fluid wedge channels converge in thickness in the direction of the rotating element's motion. The rotating element's motion applies viscous drag forces to the fluid in the converging wedge channels. This results in increases in fluid pressure, especially near the trailing end of the wedge channels. If the rotating element moves toward the non-rotating element, the wedge channel's convergence angle increases causing the fluid pressure along the wedge channel to increase. If the rotating element moves away, the pressure rise along the wedge channel decreases. Thus, the fluid in the wedge channels exerts restoring forces on the rotating element that vary with and stabilize running clearances and prevent contact between the rotating and non-rotating elements of the bearing. Flexing and sliding of the foils causes coulomb damping of any orbital motion of the rotating element of the bearing.\nAt low rotor speeds, the rotating element of the bearing is in physical contact with the fluid foil element(s) of the bearing, resulting in bearing wear. Only when the rotor speed is above what is termed the lift-off/touch-down speed do the fluid dynamic forces generated in the wedge channels assure a running gap between the rotating and non-rotating elements.\nConventional, compliant foil hydrodynamic fluid film bearings have relied on backing springs and the spring properties of the elastically bent fluid foil to preload the fluid foil against the relatively moveable rotating element (rotor) to control foil nesting/position and to establish foil dynamic stability. This preload force significantly reduces the rotor speed at which the wedge channel effect is strong enough to lift the rotating bearing element out of physical contact with the non-rotating elements of the bearing. This preload force and the high lift-off/touch-down speed result in significant bearing wear each time the rotor is started or stopped. These bearings usually have high starting torques.\nConventional, compliant foil hydrodynamic fluid film bearings achieve spring properties for their spring foils before assembly into the bearing cartridge by imparting a permanently deformed series of, for example, convoluted spring portions. Therefore, the foil pack is relatively thick and has fairly large tolerances for foil stack height. The maximum rotor deflections that these bearings permit is fairly large (loose compliance control).\nMost conventional compliant foil hydrodynamic fluid film bearings use a plurality of fluid foils and a plurality of spring foils. This results in reduced load capacity and higher bearing costs when compared to bearings with a single spring foil and a single fluid foil.\nIf a single fluid foil were to be used in lieu of a plurality of fluid foils, and that foil were formed as a flat, constant thickness sheet, which is subsequently formed to fit into the cartridge, the foil would not follow the contours of the inside of the cartridge near the foil ends. This is due to the fact that no bending moment can be applied to the ends of the foil and the foil's bending stiffness is constant throughout its length.\nIf a single fluid foil is utilized and it is attached to the cartridge at one end, the rotor can only rotate in one direction without risking foil to rotor lock-up and foil failure due to ever increasing rotor to foil friction and foil tension.\nPreviously, many types of compliant foil hydrodynamic fluid film radial bearings have been used to support, position and control the resonant motion of high speed rotating assemblies. These bearings have utilized cartridges with round apertures. Most bearings have used from five to nine fluid foils and from five to nine spring foils arrayed around the inside of the cartridge aperture. Many of the conventional, prior art bearings have the fluid foil welded to the spring foil, or when using a plurality of foils, attach the foils to the cartridge by bending the foils around bars or welding them to bars and then inserting the bars into wide slots provided for foil retention in the cartridge aperture.\nA search of the prior-art has not disclosed any directly pertinent patents to the instant invention. Specifically, no art has been found that utilizes: 1) contoured cartridge apertures, 2) flat-formed spring foils with cantilever beams which stand erect when the foil is bent or formed for insertion into the cartridge, and 3) circumferential preloading of foils to force the foils to press against the cartridge aperture. However, the following U.S. patents are considered related:\n______________________________________ U.S. PAT. NO. INVENTOR DATE ISSUED ______________________________________ 3,366,427 Silver Jan. 30, 1968 3,795,427 Licht Mar. 05, 1974 3,809,443 Cherubim May 07, 1974 3,957,317 Silver May 18, 1976 4,170,389 Eshel Oct. 09, 1979 4,178,046 Silver Dec. 11, 1979 4,213,657 Gray July 22, 1980 4,223,958 Gray Sep. 23, 1980 4,229,054 Miller Oct. 21, 1980 4,262,975 Heshmat Apr. 21, 1981 4,277,113 Heshmat July 07, 1981 4,415,281 Agrawal Nov. 15, 1983 4,552,466 Warren Nov. 12, 1985 4,743,126 Soum May 10, 1988 4,961,122 Sakai Oct. 02, 1990 5,129,739 Asai July 14, 1992 5,228,785 Saville July 20, 1993 ______________________________________\nSilver (U.S. Pat. No. 3,366,427) discloses a bead 22, wide slots 18 and 19, rod 23 as elements for attaching foils to the cartridge while in the instant disclosure the foils are retained by pressing against the circumferential preload bar or are slid into electro discharge machined EDM fabricated slots. Silver also discloses ramps which are achieved by overlapping a plurality of foils, while in the instant invention there is no overlapping of foils and the lobes or ramps are achieved by contouring the cartridge aperture.\nLicht (U.S. Pat. No. 3,795,427) discloses an overlapping single foil 61 or a plurality of overlapping foils 63, 64, and 65 which are secured to the cartridge by a retaining tab 78 at the end of the cartridge. In the instant invention, non-overlapping foils are secured by the circumferential preload bar along a line parallel to the cartridge axis.\nCherubim (U.S. Pat. No. 3,809,443) discloses a single spring foil element 23 and a single fluid foil element 25. But these foils are bonded together at one end, while in the instant invention the spring and fluid foils are not bonded. Cherubim discloses a spring foil that is formed into resilient surface elevations 27, while in the instant invention the spring foil is formed flat. Cherubim discloses a suitable restraining means (pin) 95 to secure the spring and fluid foils to the cartridge without circumferential preloading the foils, while in the instant invention a circumferential preload bar is used to preload the foils.\nSilver (U.S. Pat. No. 3,957,317) discloses overlapping foils having herringbone or chevron shaped trailing edges that are attached to the cartridge by bars set in wide milled slots. These foils are preformed with radii that are greater than that of the round cartridge aperture. The intentional mismatch of the foil/cartridge radii causes the plurality of foils to be preloaded against the rotor. In the instant invention, the single fluid foil and single spring foil are preloaded away from the rotor and towards the contoured cartridge aperture.\nEshel (U.S. Pat. No. 4,710,389) discloses a ribbon-shaped thin steel foil member 41 which is stretched over a body of pressurized material which functions as a support media for fluid foils. He does not disclose a spring foil as a support for a fluid foil.\nSilver (U.S. Pat. No. 4,178,046) discloses foils attached to the cartridge by mounting bars 14 and 44 and supported by springs with milled or chemically etched ridges 54 and grooves 56. In the instant invention, no bars are bonded to the foil or inserted into cartridge grooves nor are the spring foils formed to make ridges or grooves.\nGray (U.S. Pat. No. 4,213,657) discloses a spring comprised of tubular elements 22 and cage elements 24, while in the instant invention flat formed spring foils are utilized.\nGray (U.S. Pat. No. 4,223,958) discloses a spring foil having raised projections 42 fabricated in previously flat metal by the application of deforming forces. These projections function as cantilever beams. But in the instant invention, the cantilever beams are created in flat metal by chemical etch techniques, then they \"stand proud\" and function as springs only when the foil is bent or formed for insertion into the cartridge.\nMiller (U.S. Pat. No. 4,229,054) discloses spacer blocks 25, the ends 23 of which are attached to the bore wall of the sleeve 15 by welding or other suitable means to prevent rotation of the bearing modules 13 relative to the sleeve 15. This welding or other suitable means presumably would also function to prevent axial translation of the bearing foil modules out of the sleeve. The bearing foil modules are also shown to reside in and be captured by the walls of a circumferential channel in the sleeve, presumable to also provide constraint against the bearing foil modules translating axially out of the sleeve. Miller does not disclose any element functioning as a circumferential preload bar, nor does he disclose any element functioning as anti-translation tabs or anti-translation grooves. Miller discloses a bi-directional journal bearing utilizing three bearing foil modules, each being 120 degrees long in arc length. But he does not identify any means to generate the circumferentially compressive stresses within the foils, between the foil attachment points to the sleeve and any foil to rotor contact points which would be required for bi-directional operation of single fluid foil bearings without tightening the fluid foil about the rotor to the point of foil failure. Miller does not utilize a contoured cartridge aperture, flat formed spring foils, or spring foil spring rates that vary along the arced length of the foil.\nHeshmat (U.S. Pat. No. 4,262,975) discloses a compliant support element 54 which has a spring rate that varies along the length of the element due to variations in the spacing of ridges or bumps 56 in the element. In the instant invention, the spring rates vary due to variations in cantilever beam width and length, not variations in spacing of ridges or bumps. Heshmat also discloses a bearing pad assembly 40 which creates a ramped thickness change in the fluid foil. In the instant invention, there is a foil thickness variation, but it is in the spring foil, not the fluid foil, and it is achieved by varying the length of the cantilever beams.\nHeshmat (U.S. Pat. No. 4,277,113) discloses a spring foil with elevation 52 that has a non-linear spring rate that increase with deflection. But in the instant invention, the non-linear spring rate is accomplished by wrapping the cantilever beam around the rotor radius so as to change the free length of the cantilever beam. This is a totally different means of achieving a non-linear spring than that utilized by Heshmat.\nAgrawal (U.S. Pat. No. 4,415,281) discloses a convoluted spring foil 45 that has a non-linear spring rate that varies with deflection. This is accomplished by forming the foil in a plurality of sinusoidal shapes of varying height and spacing. Again, this method is different than in the instant invention.\nWarren (U.S. Pat. No. 4,552,466) discloses foils with tabs 24 used for attachment to the cartridge. None of the disclosed tab configurations or attachment methods utilize electro discharge machine (EDM) technology.\nSoum (U.S. Pat. No. 4,743,126) discloses a foil having side margins 3 that attach the foil to the end of the cartridge. In the instant invention, the foils are restrained against rotation in the cartridge by tabs oriented 90 degrees relative to those disclosed by Soum.\nSakai (U.S. Pat. No. 4,961,122) discloses a herringbone grooved bearing (the invention) and V-grooved bearings (the prior art) wherein the bearing has either grooves 160 or raised portions 260 in the shaft 4 and similar groove or raised portion patterns in the inner bore of the bearing housing 6 to induce local increases in the process fluid pressure when the shaft is rotated and displaced radially away from the centerline of the housing. The V grooves in the rotor and housing disclosed by Sakai are unrelated to the variable width V-shaped grooves used in the outer surface of the fluid foil of the instant invention.\nAsai (U.S. Pat. No. 5,129,739) discloses a bearing with dynamic pressure grooves formed in the thermoplastic resin inner surface of a rigid metal cylinder 21 which functions as a fluid dynamic surface. This cylinder is not resiliently mounted to the bearing housing. In the fourth embodiment of the instant invention, the fluid sleeve does not have pressure grooves in its inner surface, is thin and flexible, and is resiliently mounted to the bearing cartridge aperture.\nSaville (U.S. Pat. No. 5,228,785) discloses overlapping foils supported by convoluted spring foils having variable convolute height and variable convolute spacing to achieve non linear spring rates and ramps. In the instant invention, non-linear spring rates are achieved with variable cantilever beam width and length, and ramps are achieved by contouring the cartridge aperture and with variable cantilever beam length."}
{"text": "The disclosure of Japanese Patent Application Nos. 2001-328425 filed on Oct. 25, and 2002-232195 filed on Aug. 9, 2002 including the specification, drawings and abstract are incorporated herein by reference in their entirety.\n1. Field of the Invention\nThe invention relates to a motor-driven type fuel pump for a vehicle to be accommodated in a fuel tank for a vehicle, and more particularly to a motor-type fuel pump for a vehicle to reliably prevent entry of a fuel compressed by a compression stroke of a pump into an intake side as vapor, and further to allow a single hole to serve both as positioning pin for executing positioning when assembling a housing and a cover and as a pressure relief hole for releasing pressure in a pump chamber.\n2. Description of Related Art\nRecently, as a fuel pump for supplying fuel to a vehicle engine, an in-tank type motor-driven fuel pump which is installed in a fuel tank is used. Among them, a fuel pump suspending from a flange member secured to an opening of an upper wall of the fuel tank is widely used. Further, a unit-type fuel pump incorporating a filter or the like thereinto has also been used.\nIn such fuel pumps, a pump body portion is structured by combining a housing 42 formed with a pump chamber 41 therein and a cover 43 abutting against a lower surface of the housing 42 so as to cover the pump chamber 41, as shown for example in a sectional view of the pump portion in FIG. 9. The pump chamber 41 is provided, as a pump member, with an inner rotor 44 of a trochoid-gear type or the like. The inner rotor 44 is rotated by a rotation axis 46 of a motor 45, and a fuel in the fuel tank is sucked through an intake port 47 formed on the cover 43 and discharged through a discharge hole 48 into a motor chamber 50. Next, the fuel passing through the motor chamber 50 is pressurized, and supplied through an exhaust hole, not shown, on the upper portion of the pump body to a fuel injection system and the like.\nIn a fuel pump like this, particularly when the fuel temperature is increased, there are some cases where the fuel pressurized inside the pump generates vapor from the inside of the pump and deteriorates pump performance. Therefore, as shown in FIG. 9B, a pressure relief hole 54 communicated with the outside of the pump is formed piercing the cover 43 at a position between the intake port 47 and the discharge hole 48 at an end portion of a place at which a high-pressure area of the pump is formed.\nFIGS. 10A and 10B show this state in more detail. Particularly, as shown in FIG. 10B, which is a sectional view taken along line XBxe2x80x94XB portion in FIG. 10A, an internal tooth gear-like pump chamber outer peripheral surface is formed on an inner peripheral surface 56 of the pump chamber portion in an outer rotor 60 rotatably disposed in the housing 42. Inside the internal tooth gear-like pump chamber outer peripheral surface, rotatably provided is an external tooth gear-like inner rotor 44, capable of meshing with the pump chamber outer peripheral surface in a manner of the trochoid-gear type, and having a small diameter and a smaller number of teeth than the aforementioned internal teeth. Then, the inner rotor 44 is rotated by the rotation axis 46 driven by the motor 45 as above, and pump action is carried out.\nAt this time, in accordance with rotation of the inner rotor 44 in a direction shown in arrow R in the figure, the fuel sucked through the intake port 47 is sucked via an intake groove 57 formed on the upper surface of the cover 43 into a chamber portion at an intake stroke side in the pump chamber. After that, the chamber at the intake stroke side is closed and the fuel is compressed, the chamber is communicated with a discharge groove 58 formed at a discharge port 48 side, and the fuel is discharged through the discharge port 48 to the inside of the motor body.\nWhen the inner rotor 44 is further rotated, the external teeth of the inner rotor 44 mesh with the internal teeth formed on the inner surface of the outer rotor 60, thereby forming a high-pressure closed chamber portion P in the figure. After that, when the high-pressure closed chamber is communicated with the intake groove 57 caused by the rotation of the inner rotor 44, a high-pressure fuel enters a low-pressure fuel through the intake hole 47. Since the pressure is released, liquid inside carries out vacuum boiling and vaporizes. Therefore, vapor contaminates the fuel at the intake side in the pump chamber, causing a vapor lock state in which the fuel is unable to be compressed sufficiently. This may sometimes deteriorate pump performance remarkably.\nAs a countermeasure, for example, as shown in FIGS. 9A, 9B, 10A and 10B, the pressure relief hole 54 is formed at a position where a high-pressure area of the pump is formed. The pressure relief hole 54, for example, as shown in FIG. 10B, is provided on the cover 43 between a front edge portion of the intake groove 57 and a rear edge portion of the discharge groove so as to pierce the cover toward below the pump. Owing to this, the high-pressure fuel is introduced from a portion of the highest pressure shown as point P in the figure to the pressure relief hole 54, through a narrow gap between opposite wall surfaces of the cover 43 and the housing 42. Then, the fuel is released from the pressure relief hole 54 into the fuel tank.\nA problem of generation of vapor lock due to entry of the high-pressure fuel at the discharge side into the intake side in the fuel pump as above occurs not only in a trochoid-gear type positive-displacement fuel pump as shown in the figure but also in other positive-displacement fuel pumps. Further, as well as in positive-displacement fuel pumps, similar problems occur in various non-positive-displacement fuel pumps of wesco type or the like.\nIn the meantime, when assembling the housing 42 and the cover 43 as above, in order to easily perform a precise relative positioning between the intake port 47 portion formed on the cover 43 and the discharge hole 48 portion formed on the housing 42 or the like, a positioning pin 52 is inserted into a positioning pin insertion hole 59 at the housing 42 side and fixed in advance so as to protrude toward the cover 43 side, which is then inserted into the positioning pin insertion hole 59 formed in the cover 43. Alternatively, the positioning pin 52 is inserted into the positioning pin insertion hole 59 of the cover 43 and fixed so as to protrude toward the housing 42 side, which is then inserted into the positioning pin insertion hole 59 formed on the housing 42.\nAccordingly, when assembling this fuel pump, for example, in a case where the positioning pin 52 is fixed at the housing 42 side in advance, a tip of the positioning pin 52 is fitted into the pin insertion hole 53 of the cover 43 so as to combine the both in a state where a pump member 44 is disposed in the pump chamber 41 of the housing 42. The combined one is incorporated into a tip of the rotation axis 46 of the motor 45, and the tip and the rotation axis are bonded by an external casing 55, so as to provide an integrated fuel pump unit as a whole. Note that, an example of the positioning pin 52 to be used includes a C-type pin in cross section or the like.\nIn a fuel pump like this, in order to prevent generation of vapor lock caused by entry of the fuel across the boundary area from the high-pressure discharge side to the low-pressure intake side in the pump, the pressure relief hole 54 for communicating the boundary area with the outside of the pump was provided. However, as the pressure relief hole is provided at one point of the area, namely in a point-like manner, pressure relief was not always carried out sufficiently. Thus, effects of preventing generation of vapor lock is not sufficient.\nFurther, in a fuel pump as above, it was necessary to form the positioning pin insertion hole 53 in the cover 43 through which the positioning pin 52 passes, and further the pressure relief hole 54 as above. Therefore, since it was necessary to form these holes separately in the cover 43 or the like of the pump, man-hours are increased and the pump becomes expensive.\nTherefore, it is a first object of the invention to prevent with reliability generation of vapor lock caused by entry of high-pressure fuel at a discharge side of a pump chamber into an intake side thereof. It is a second object of the invention to eliminate the necessity of forming a positioning pin press fitting hole and a pressure relief hole separately, allow a single hole to serve as both holes. Further, it is a third object of the invention to achieve the first and second objects by using the same means.\nIn order to achieve the first object, a motor-type fuel pump for a vehicle according to a first aspect of the invention is provided with a pressure relief portion extended along a boundary area portion between a high-pressure portion and a low-pressure portion, which be communicated with a pressure chamber so as to reduce a pressure in the pressure chamber\nNote that, in a case of a trochoid-gear type pump as shown in the figure, the xe2x80x9chigh-pressure portionxe2x80x9d of the pump refers to a high-pressure portion in which a sucked fuel is compressed due to reduced volume of a pump chamber in accordance with rotation of an inner rotor and an outer rotor, and xe2x80x9clow-pressure portionxe2x80x9d refers to a low-pressure portion in which the fuel is sucked and secured due to increased volume of the pump chamber in accordance with rotation of the inner rotor and the outer rotor. Note that, a problem like this occurs not only in a trochoid-gear type pump like this, but also in other positive-displacement fuel pumps. Further, it also occurs in various non-positive-displacement fuel pumps of a wesco type or the like as well as postivie-displacement fuel pumps. In those pumps, there are also a xe2x80x9chigh-pressure portionxe2x80x9d which is a compression side, and a xe2x80x9clow-pressure portionxe2x80x9d which is an intake side. The problem like this can be solved by using the aforementioned means for solving the problem.\nFurther, a fuel pump according to the invention may have a structure in which, at a boundary area portion between the high-pressure portion and the low-pressure portion of the pump, a groove or a slit is formed extending along the boundary area portion, so that the groove or the slit communicates the chamber with the outside of the fuel pump.\nFurther, in order to achieve the second object, a motor-type fuel pump for a vehicle according to a second aspect of the invention, is provided with a pressure relief hole which both positions a housing that accommodates the pump therein and a cover formed with a fuel intake hole with a positioning pin and releases pressure in a pump chamber. The motor-type fuel pump is structured such that the positioning pin is inserted into the aforementioned pressure relief hole and a pressure relief passage can be secured in the aforementioned pressure relief hole when inserting the positioning pin.\nFurther, in order to achieve the aforementioned first and second objects by the same method, a fuel pump according to the second aspect of the present invention may be structured such that the aforementioned pressure relief hole may be communicated with the boundary area portion between the high-pressure portion and the low-pressure portion by means of a opening, such as a pressure relief groove, extending along the boundary area portion."}
{"text": "1. Field of the Invention\nThe present invention relates to artificial fusion implants to be placed into the intervertebral space left remaining after the removal of a damaged spinal disc and specifically to the apparatus for and method of, inserting the implants.\n2. Description of the Prior Art\nFor the purpose of achieving long term stability to a segment of injured spine, a fusion (the joining together of two or more bones via a continuous bridge of incorporated bone) may be performed. Well-known to those skilled in such art is the interbody fusion wherein the disc is partially excised and bone placed within that space previously occupied by that disc material (between adjacent vertebrae) for the purpose of restoring a more normal spatial relationship, and to provide for stability; short term by mechanical support, and long term by the permanent cross bonding of bone from vertebra to vertebra. For fusion to occur within the disc space, it is necessary to prepare the vertebrae to be fused by breaking through, or cutting into, the hardened outside plates of bone (the endplates) to allow the interposed bone graft to come into direct contact with the more vascular cancellous (spongy) bone, and to thereby trick the body into attempting to heal this induced, but controlled, \"fracturing\" by both bone production and the healing of the grafts to both opposed vertebral surfaces such that they become one continuous segment of bone.\nThe purpose of the present invention is to provide an implant, and the apparatus and method of inserting the implant within the intervertebral space left after the removal of the disc material and permanently eliminate all motion at that location. To do so, the device of the present invention is space occupying within the disc interspace, rigid, self-stabilizing to resist dislodgement, stabilizing to the adjacent spinal vertebrae to eliminate local motion, and able to intrinsically participate in a vertebra to vertebra bony fusion so as to assure the permanency of the result.\nAt present, following the removal of a damaged disc, either bone or nothing is placed into the remaining space. Placing nothing into this space allows the space to collapse which may result in damage to the nerves; or the space may fill with scar tissue and eventually lead to a reherniation. The use of bone to fill the space is less than optimal in that bone obtained from the patient requires additional surgery and is of limited availability in its most useful form, and if obtained elsewhere, lacks living bone cells, carries a significant risk of infection, and is also limited in supply as it is usually obtained from accident victims. Furthermore, regardless of the source of the bone, it is only marginal structurally and lacks a means to either stabilize itself against dislodgement, or to stabilize the adjacent vertebrae.\na. Prior Art Implants\nThere have been an extensive number of attempts to develop an acceptable disc prosthesis (an artificial disc). Such devices by design would be used to replace a damaged disc and seek to restore the height of the interspace and to restore the normal motion of that spinal joint. No such device has been found that is medically acceptable. This group of prosthetic or artificial disc replacements, seeking to preserve spinal motion and so are different from the present invention, would include:\nU.S. Pat. No. 3,867,728 to STUBSTAD--describing a flexible disc implant.\nU.S. Pat. No. 4,349,921 to KUNTZ--describing a flexible disc replacement with file-like surface projections to discourage device dislocation.\nU.S. Pat. No. 4,309,777 to PATIL--describing a motion preserving implant with spiked outer surfaces to resist dislocation and containing a series of springs to urge the vertebrae away from each other.\nU.S. Pat. No. 3,875,595 to FRONING--describing a motion preserving bladder-like disc replacement with two opposed stud-like projections to resist dislocation.\nU.S. Pat. No. 2,372,622 to FASSIO (France)--describing a motion preserving implant comprising complimentary opposed convex and concave surfaces.\nIn summary, these devices resemble the present invention only in that they are placed within the intervertebral space following the removal of a damaged disc. In that they seek to preserve spinal motion, they are diametrically different from the present invention which seeks to permanently eliminate all motion at that spinal segment.\nA second related area of prior art includes those devices utilized to replace essentially wholly removed vertebrae. Such removal is generally necessitated by extensive vertebral fractures, or tumors, and is not associated with the treatment of disc disease. While the present invention is to be placed within the disc space, these other vertebral devices cannot be placed within the disc space as at least one vertebra has already been removed such that there no longer remains a \"disc space\". Furthermore, these devices are limited in that they seek to perform as temporary structural members mechanically replacing the removed vertebrae (not a removed disc), and do not intrinsically participate in supplying osteogenic material to achieve cross vertebrae bony fusion. Therefore, unlike the present invention which provides for a source of osteogenesis, use of this group of devices must be accompanied by a further surgery consisting of a bone fusion procedure utilizing conventional technique. This group consisting of vertebral struts rather than disc replacements would include the following:\nU.S. Pat. No. 4,553,273 to WU--describing a turnbuckle-like vertebral strut.\nU.S. Pat. No. 4,401,112 to REZAIAN--describing a turnbuckle-like vertebral strut with the addition of a long stabilizing staple that spans the missing vertebral body.\nU.S. Pat. No. 4,554,914 to KAPP--describing a large distractible spike that elongates with a screw mechanism to span the gap left by the removal of an entire vertebra and to serve as an anchor for acrylic cement which is then used to replace the missing bone (vertebrae).\nU.S. Pat. No. 4,636,217 to OGILVIE--describing a vertebral strut mechanism that can be implanted after at least one vertebrae has been removed and consists of a mechanism for causing the engagement of screws into the vertebrae above and the vertebrae below the one removed.\nIn summary, this second group of devices differs from the present invention in that they are vertebral replacements struts, do not intrinsically participate in the bony fusion, can only be inserted in the limited circumstances where an entire vertebra has been removed from the anterior approach, and are not designed for, or intended to be used for the treatment of disc disease.\nA third area of prior art related to the present invention includes all devices designed to be applied to one of the surfaces of the spine. Such devices include all types of plates, struts, and rods which are attached by hooks, wires and screws. These devices differ significantly from the present invention in that they are not inserted within the disc space and furthermore do not intrinsically participate in supplying osteogenic material for the fusion.\nTherefore, where permanent spinal immobilization is desired, an additional surgery, consisting of a spinal fusion performed by conventional means or the use of supplemental methylmethacrylate cement is required. Such devices applied to the spine, but not within the disc space, would include the following:\nU.S. Pat. No. 4,604,995 to STEPHENS--describing a \"U\" shaped metal rod attached to the posterior elements of the spine with wires to stabilize the spine over a large number of segments.\nU.S. Pat. No. 2,677,369 to KNOWLES--describing a metal column device to be placed posteriorly along the lumbar spine to be held in position by its shape alone and to block pressure across the posterior portions of the spinal column by locking the spine in full flexion thereby shifting the maximum weight back onto the patient's own disc.\nOther devices are simply variations on the use of rods (e.g. Harrington, Luque, Cotrel-Dubosset, Zielke), wires or cables (Dwyer), plates and screws (Steffee), or struts (Dunn, Knowles)\nIn summary, none of these devices are designed to be nor can be used within the disc space. Moreover, these devices do not replace a damaged disc, and do not intrinsically participate in the generation of a bony fusion.\nAnother area of related prior art to be considered is that of devices designed to be placed within the vertebral interspace following the removal of a damaged disc, and seeking to eliminate further motion at that location.\nSuch a device is contained in U.S. Pat. No. 4,501,269 issued to BAGBY which describes an implantable device and limited instrumentation. The method employed is as. follows: a hole is bored transversely across the joint and a hollow metal basket of larger diameter than the hole is then pounded into the hole and then the hollow metal basket is filled with the bone debris generated by the drilling.\nWhile the present invention (device, instrumentation, and method) may appear to bear some superficial resemblance to the BAGBY invention, it is minimal, while the differences are many fold and highly significant. These differences include the following:\n1. Safety--The present invention provides for a system of completely guarded instrumentation so that all contiguous vital structures (e.g. large blood vessels, neural structures) are absolutely protected. The instrumentation of the present invention also makes overpenetration by the drill impossible. Such overpenetration in the cervical spine, for example, would result in the total paralysis or death of the patient. In the thoracic spine, the result would be complete paraplegia. In the lumbar spine, the result would be paraplegia or a life-threatening perforation of the aorta, vena cava, or iliac vessels. PA1 2. Broad applicability--The BAGBY device can only be inserted from the front of the vertebral column, however, in contrast, the present invention can be utilized in the cervical, thoracic, and lumbar spine, and can be inserted from behind (posteriorly) in the lumbar spine. This is of great importance in that the purpose of these devices is in the treatment of disc disease and probably greater than 99 percent of all lumbar operations for the treatment of disc disease are performed from behind where the present invention can easily be utilized, but the BAGBY device, as per BAGBY'S description, cannot. PA1 3. Disc removal--The BAGBY invention requires the complete removal of the disc prior to the drilling step, whereas the present invention eliminates the laborious separate process of disc removal and efficiently removes the disc and prepares the vertebral end plates in a single step. PA1 4. Time required--The present invention saves time over the BAGBY invention since time is not wasted laboring to remove the disc prior to initiating the fusion. Also, with the present invention the procedure is performed through a system of guarded instrumentation, time is not wasted constantly placing and replacing various soft tissue retractors throughout the procedure. PA1 5. Implant stability--Dislodgement of the implant would be a major source of device failure (an unsuccessful clinical result), and might result in patient paralysis or even death. As discussed, the BAGBY device lacks any specific means of achieving stability and since it is pounded in against resistance to achieve vertebral distraction, and is susceptible to forceful dislodgement by the tendency of the two distracted vertebrae, to return to their original positions squeezing out the device. The present invention, however, is screwed into place. As there is no unscrewing force present between the vertebrae, compression alone cannot dislodge the implant. The implant is inherently stable by its design. Furthermore, the threads of the present invention are highly specialized in that they are periodically interrupted so that the tail ends of each of the tabs so formed are blunted and twisted so as to resist accidental unscrewing. The removal of an implant with such \"locking threads\" requires the use of a special extractor included within the instrumentation. The stability of the present invention is still further enhanced, again in contradistinction to the BAGBY device, by the presence of a \"bone ingrowth\" surface texturing, which both increases the friction of the fit and allows for the direct growth of the vertebral bone into the casing of the implant itself. PA1 6. Spinal stability--The present invention is not only self-stabilizing, it also provides stability to the adjacent vertebrae in at least three ways that the BAGBY device cannot. First, the BAGBY device is placed transversely across the joint in the center, leaving both vertebrae free to rock back and forth over this round barrel shaped axis, much like a board over a barrel, being used for a seesaw. PA1 7. The collapse of the interspace--While both the present invention and the BAGBY device can be fabricated to withstand the compression forces within the interspace, the interspace may nevertheless collapse under the superincumbent body weight as the implant settles into the vertebral bone. This is related to the load per unit area. Again the present invention is superior to the BAGBY device in at least four ways. PA1 8. Bone ingrowth surface texturing--The present invention has a surface treatment of known and conventional technology to induce the growth of bone from the vertebrae directly into the casing material of the implant itself. The BAGBY device has no similar feature. PA1 9. Fusion mass--The BAGBY invention calls for removing the disc and then drilling a hole between the adjacent vertebrae. The bony debris so generated is then put into the device. The present invention takes a core of pure bone producing marrow from the iliac crest, and then by use of a special press, forcibly injects the implant device with an extremely dense compressed core of that osteogenic material until the material itself virtually extrudes from every cell of the implant. PA1 10. The probability of achieving fusion--The fusion rate within the spine is known to be related directly to the amount of exposed vascular bone bed area, the quality and quantity of the fusion mass available, and the extent of the stabilization obtained with all other factors being half constant. It would then be anticipated, that the fusion rate would be superior with the present invention as compared to the BAGBY device, because of optimal implant stability (#5), optimal spinal stability (#6), bone ingrowth surface treatment (#8), superior fusion mass (#9), and the greater exposed vertebral bony surface area (#7). PA1 1. Failure to protect the surrounding tissues throughout the procedure, specifically, prior to drilling and until after the insertion of the graft; PA1 2. Failure to contain the debris, bony and cartilaginous, generated during the procedure; PA1 3. Failure to optimize the contact of the cylindrical drill hole and bone graft, the mismatch in their diameters resulting in incongruence of fit; PA1 4. Failure to determine the optimal drill size prior to drilling; PA1 5. Failure to determine the optimal amount of distraction prior to drilling; PA1 6. Inability to optimize the amount of distraction so as to restore the normal spatial relationships between adjacent vertebrae; PA1 7. Inability to create sufficient working space within the spinal canal (between the nerve roots and the dural sac) to make the procedure safe; PA1 8. Absent a foot plate on the drill guide, as necessitated by the close tolerances posteriorly, the inability to reliably insure that the drilling is parallel to the vertebral endplates; PA1 9. The inability to insure equal bone removal from the opposed vertebral surfaces; and PA1 10. The inability to determine within the spinal canal, the proper side by side positioning for dual drill holes. PA1 1. The present invention is safer by providing protection of the surrounding tissues. An Outer Sleeve places all of the delicate soft tissue structures, nerves, blood vessels, and organs outside of the path of the various sharp surgical instruments and the implant. Further, it is an improvement upon hand held retractors in that it occupies the least possible amount of area, avoids the stretching associated with manual retraction, provides for the retraction and shielding of the surrounding tissues in all directions circumferentially and simultaneously, and it does so exclusively with smooth, curved surfaces. PA1 2. The present invention is safer by providing protection against the danger of instrument or implant overpenetration. PA1 3. The present invention is safer as the surgical site and wound are protected from the debris generated during the procedure. PA1 4. The present invention is safer because the method provides for absolute protection to the soft tissues directly and from indirect injury by overpenetration. It makes safe the use of power instrumentation which is both more effective and efficient. PA1 5. The present invention maintains the vertebrae to be fused rigid throughout the procedure. PA1 6. The present invention holds the vertebrae to be fused aligned throughout the procedure. PA1 7. The present invention holds the vertebrae to be fused distracted throughout the procedure. PA1 8. The present invention assures that all instruments introduced through the Outer Sleeve are coaxial and equally centered through the disc space and parallel the endplates. PA1 9. The present invention facilitates the implant insertion by countering the high compressive forces tending to collapse the interspace, which if left unchecked would resist the introduction and advancement of the implant and make stripping more likely. PA1 10. The present invention extends the range and use of the procedure and similarly the interbody spinal implant itself by making the procedure safe throughout the spine. PA1 11. The present invention increases the ability to use a specifically sized implant. PA1 12. In the present invention the end of all the penetrating instrumentation is blunt faced. PA1 13. In the present invention all of the instruments have been stopped at a predetermined depth to avoid overpenetration. PA1 14. The design of the Outer Sleeve in the present invention conforms to the spacial limitations of the specific surgical site. PA1 15. The design and use of a second or Inner Sleeve in the present invention allows for the difference in size between the inside diameter of the Outer Sleeve, and the outside diameter of the drill itself. This difference being necessary to accommodate the sum of the distraction to be produced, and the depth of the circumferential threading present of the implant. PA1 16. In the present invention a specially designed drill bit with a central shaft recess allows for the safe collection of the drilling products, which can then be removed without disturbing the Outer Sleeve by removing the drill bit and Inner Sleeve as a single unit. PA1 17. In the present invention a specially designed trephine is provided for preparing cylindrical holes across the disc space and into two adjacent vertebrae, and for harvesting of the bone graft by removing a core of bone slightly smaller in diameter than the internal diameter of the implant cavity itself, however of a greater length. PA1 18. In the present invention a specially designed press for forcefully compressing and injecting the long core of autogenous bone into the implant, such that it extrudes through the implant itself. PA1 19. In the present invention a specially designed driver extractor, which attaches to the implant and allows the implant to be either inserted or removed without itself dissociating from the implant, except by the deliberate disengagement of the operator. PA1 20. In the present invention predistraction increases the working space. PA1 21. The Distractor in the present invention is self-orienting acting as a directional finder. PA1 22. The Distractor in the present invention is self-centralizing between the opposed vertebral surfaces acting as a centering post for the subsequent bone removal. PA1 23. In the present invention predistraction assures the equal removal of bone from the adjacent vertebral surfaces. PA1 24. In the present invention predistraction assures the exact congruence between the hole drilled and the device. PA1 25. In the present invention predistraction assures that the drilling is parallel to the vertebral endplates. PA1 26. In the present invention predistraction allows for the determination of the optimal distraction prior to drilling. PA1 27. In the present invention predistraction allows for the verification of the correct prosthesis size prior to drilling. PA1 28. In the present invention predistraction facilitates device insertion by relieving the compressive loads across the interspace which would resist implantation. PA1 29. In the present invention predistraction decreases the likelihood of stripping the bone during insertion. PA1 30. In the present invention predistraction provides for the side by side positioning, spacing, and parallelism required prior to the irrevocable event of drilling. PA1 31. In the present invention predistraction provides for the rigid stabilization of the vertebrae opposed to the disc space throughout the surgical procedure. PA1 32. In the present invention predistraction provides for an implant easier to insert as the compressive loads of the opposed vertebrae are held in check so that the device itself need not drive the vertebrae apart to be inserted. PA1 33. In the present invention predistraction allows for the insertion of a more effective implant as more of the implant can be dedicated to its intended purpose and be full diameter, whereas without the benefit of predistraction and the ability to maintain the same, a significant portion of the forward end of the implant would need to be dedicated to the purpose of separating the opposing vertebrae. PA1 34. The present invention allows for the use of an implant with a sharper thread or surface projections as there is no danger to the surrounding tissues. PA1 35. The present invention allows for the implant to be fully preloaded as provided to the surgeon, or for the surgeon to load it with the material of his choice at the time of surgery. PA1 36. The present invention allows for the loading of a spinal implant outside of the spinal canal and prior to implantation. PA1 37. The present invention provides for Distractors for the restoration of the physiologic amount of lordosis/kyphosis at any given level of the spine. PA1 38. The Distractors of the present invention provide for the restoration of sagittal spinal alignment and the segmental correction of scoliosis. PA1 39. The present invention provides Extended Outer Sleeves which are a combined Distractor and Outer Sleeve for the three dimensional segmental restoration of spinal alignment.\nThe present invention is atraumatically screwed into place while the BAGBY device, in contradistinction, is pounded into position. BAGBY describes that its implant is significantly larger in size than the hole drilled and must be pounded in. This is extremely dangerous and the pounding occurs directly over the spinal cord which is precariously vulnerable to percussive injury. Furthermore, while it is possible, for example in the lumbar spine, to insert the present invention away from the spinal cord and nerves, the BAGBY device must always be pounded directly towards the spinal cord.\nFurthermore, since the BAGBY device is pounded into a smooth hole under great resistance, and lacking any specific design features to secure it, the device is highly susceptible to forceful ejection which would result in great danger to the patient and clinical failure. The present invention, in contradistinction, is securely screwed into place, and possesses highly specialized locking threads to make accidental dislodgement impossible. Because of the proximity of the spinal cord, spinal nerves, and blood vessels, any implant dislodgement as might occur with the BAGBY device might have catastrophic consequences.\nSecondly, as the BAGBY device lacks any specific design features to resist sliding, it may actually behave as a third body allowing the translation of the vertebrae relative to the device and to each other.\nThirdly, any device can only provide stability if it remains properly, seated. The present invention is inherently stable, and therefore assures that it will stabilize the adjacent vertebrae, rather than, as with the BAGBY, the instability of the spine to be treated may cause a dislocation of the BAGBY implant, with further loss of spinal stability.\nFirst, the present invention offers considerably greater surface area to distribute the load. Secondly, while the BAGBY device is placed centrally, the present device is placed bilaterally where the bone tends to be more cortical and much stronger out towards the rim. Thirdly, the present invention supports the load achieving an \"I\" beam effect, whereas the BAGBY implant does not. Fourthly, it is not pressure alone that causes the collapse of the bone adjacent to the implant, but also bony erosion that is caused by the motion under pressure of the implant against the bone. As discussed in item 6 above, the present invention alone is highly resistant to such motion, again diminishing the likelihood of erosion and interspace collapse.\nThe last area of prior art possibly related to the present invention and therefore, to be considered related to \"bony ingrowth\", are patents that either describe methods of producing materials and or materials or devices to achieve the same. Such patents would include:\nU.S. Pat. Nos. 4,636,526 (DORMAN), No. 4,634,720 (DORMAN), No. 4,542,539 (ROWE), No. 4,405,319 (COSENTINO), No. 4,439,152 (SMALL), No. 4,168,326 (BROEMER), No. 4,535,485 (ASHMAN), No. 3,987,499 (SCHARBACH), No. 3,605,123 (HAHN), No. 4,655,777 (DUNN), No. 4,645,503 (LIN), No. 4,547,390 (ASHMAN), No. 4,608,052 (VAN KAMPEN), No. 4,698,375 (DORMAN), No. 4,661,536 (DORMAN), No. 3,952,334 (BOKROS), No. 3,905,047 (LONG), No. 4,693,721 (DUCHEYNE), No. 4,070,514 (ENTHERLY).\nHowever, while the implant of the present invention would utilize bone ingrowth technology, it would do so with conventional technology.\nb. Prior Art Instrumentations And Methods\nThe following is a history of the prior art apparatus and methods of inserting spinal implants:\nIn 1956, Ralph Cloward developed a method and instruments which he later described for preparing the anterior aspect (front) of the cervical spine, and then fusing it. Cloward surgically removed the disc to be fused across and then placed a rigid drill guide with a large foot plate and prongs down over an aligner rod and embedded said prongs into the adjacent vertebrae to maintain the alignment so as to facilitate the reaming out of the bone adjacent the disc spaces. As the large foot plate sat against the front of the spine, it also served as a fixed reference point to control the depth of drilling. The reaming left two opposed resected arcs, one each, from the opposed vertebral surfaces. The tubular drill guide, which was placed only preliminary to the drilling, was thereafter completely removed. A cylindrical bony dowel, significantly larger in diameter than the hole formed, was then pounded into the hole already drilled. Cloward's method of instrumentation was designed for, and limited to, use on the anterior aspect and in the region of the cervical spine only. The hole was midline, which would preclude its use posteriorly where the spinal cord would be in the way.\nAs the bone graft to be inserted in Cloward's method was necessarily larger in diameter than the hole drilled, the graft could not be inserted through the drill guide. This mandated the removal of the drill guide and left the graft insertion phase completely unprotected. Thus Cloward's method and instrumentation was inappropriate for posterior application.\nIn addition, the failure to provide continuous protection to the delicate neural structures from the instruments, as well as the bony and cartilaginous debris generated during the procedure, made Cloward's method inappropriate for posterior application. Also, the drill guide described by Cloward could not be placed posteriorly within the spinal canal, as the foot plate would crush the nerves. Modifying Cloward's drill guide by removing the foot plate completely, would still leave the instrument unworkable as it would then lack stability, and would not be controllable for depth of seating.\nNevertheless, Wilterberger, (Wilterberger, B. R., Abbott, K. H., \"Dowel Intervertebral Fusion as Used in Lumbar Disc Surgery,\" The Journal of Bone and Joint Surgery, Volume 39A, pg. 234-292, 1957) described the unprotected drilling of a hole from the posterior into the lumbar spine between the nerve roots and across the disc space, and then inserting a stack of button-like dowels into that space. While Wilterberger had taken the Cloward concept of circular drilling and dowel fusion and applied it to the lumbar spine from a posterior approach, he had not provided for an improved method, nor had he advanced the instrumentation so as to make that procedure sufficiently safe, and it rapidly fell into disrepute.\nCrock (Crock, H. V., \"Anterior Lumbar Interbody Fusion--Indications for its Use and notes on Surgical Technique,\" Clinical Orthopedics, Volume 165, pg. 157-163, 1981) described his technique and instrumentation for Anterior Interbody Fusion of the lumbar spine, wherein he drilled two large holes side by side across the disc space from anterior to posterior essentially unprotected and then pounded in two at least partially cylindrical grafts larger than the holes prepared.\nA review of the prior art is instructive as to a number of significant deficiencies in regard to the method and instrumentation for the performance of Interbody Spinal Fusion utilizing drilling to prepare the endplates.\nAs the great majority of spinal surgery is performed in the lumbar spine and from posteriorly, a review of the prior art reveals a number of deficiencies in regard to the spine in general, and to the posterior approach to the lumbar spine specifically. These deficiencies include the:"}
{"text": "Integrated circuits are formed on semiconductor substrates, or wafers. The wafers are then sawed into microelectronic dies (or “dice”), or semiconductor chips, with each die carrying, for example, a respective integrated circuit (e.g., a microprocessor) or a microelectromechanical system (MEMS) device (e.g., an accelerometer). In some examples, each semiconductor chip is mounted to a package or carrier substrate using either wirebonding or “flip-chip” connections. The packaged chip is then typically mounted to a circuit board, or motherboard, before being installed in a system, such as an electronic or a computing system.\nWhile lower frequency applications typically use direct conductive connections (e.g., vias and solder formations) to connect to the circuit board, higher frequency applications, such as millimeter wave devices, often use electromagnetic coupling and a waveguide structure. In order to achieve the desired performance in the higher frequency applications, the package substrate must be precisely manufactured, which typically involves manufacturing processes that are not compatible with conventional semiconductor manufacturing processes. As a result, manufacturing costs are increased.\nAccordingly, it is desirable to provide an improved method for forming a microelectronic assembly for use in high frequency applications that allows the required performance to be achieved while minimizing costs. Other desirable features and characteristics of the invention will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "1. Field of the Invention\nThis invention relates to devices for storing and displaying model airplanes, and more particularly to devices for mounting such airplanes to wails. Still more specifically, it concerns an adjustable wall-mounted bracket assembly for safely and conveniently supporting model airplanes of various sizes, weights, and configurations.\n2. Prior Art\nStoring small model airplanes, that is, those with a wingspan of up to, say, 2 feet and weighing 12 to 14 ounces poses few problems. Because of their small size and light weight, bench or shelf space can generally be found for them, and long periods of rest on their landing gear or supported horizontally on a table-top stand impose little if any structural stress. If space is at a premium, models of this class frequently are stored vertically by hanging them from wire or cord tethers secured to a wall and looped over or attached to their nose or empennage. For larger, heavier models, these storage and display methods are not satisfactory.\nShelf or tabletop space is generally not available for models with wingspans greater than 2 feet. Suspending larger, heavier models for long periods by means of hooks, rings, or straps attached to or around their nose or tail section is almost certain to result in the attachment hardware damaging the fuselage or tail component. In addition, since the typical single-point suspension device does not provide means for preventing the model from swinging and rotating, there is a likelihood that over time air currents and careless passers-by will cause surface or structural damage.\nThe availability of more powerful engines and the development of improved construction techniques and materials over the past 20 years have led to growing interest among enthusiasts in ever bigger, heavier models. \"Jumbo\" scale planes having wingspreads reaching over 7 feet and weights exceeding 30 pounds are not uncommon. Even when space for the horizontal display of such models is available, storing them for long periods on their landing gear creates flat spots on their tires and warped wings.\nWith the proliferation of models with spans in excess of three feet and weights exceeding 12 pounds, the need for a reliable storage system has become acute. The conventional response to this need is to hang it, nose upward or downward, on a pair of horizontally spaced wall-mounted supports, such as wire hooks bent to conform roughly to the contour of the trailing or leading edge of the model's wings. The model's weight is borne by the wings and rests on their trailing or leading edges at their points of contact with the hooks. Because of the great and ever growing disparity in the sizes and configurations of this category of models, most prior art mounting devices are either designed specifically for a particular model, with wing- or fuselage-supporting hardware attached to a wall-mounted backing, or offered as individual pieces of mounting hardware, such as common metal L-, or U-shaped brackets adapted to be attached to a wall at spaced intervals conforming to the configuration and size of the particular model. Both of these classes of devices suffer from a number of deficiencies. For one, neither of them is suitable for long-term storage. Most model aircraft wings are not designed to withstand substantial forces applied cordwise of the wing. Supporting a large, heavy model by its wings alone in either a nose-down or nose-up attitude over time can result in bending and warping of wing spars, ribs, and other structural components, twisting of fittings and fairings, and stressing of the aircraft skin.\nAnother deficiency concerns the limited utility of the prior art wall-support systems. The devices in which the brackets are joined or mounted to a backing are designed for use with a specific airplane model. Generally they either cannot be modified, or lend themselves only with difficulty to reconfiguration for use with another model. From a manufacturing standpoint, these devices create a requirement for diversity in product sizes and shapes which translates into unacceptably high production and inventory costs and stocking nightmares. Similarly for distributors and retailers, maintaining an adequate inventory to match even the most popular models poses a difficult and expensive challenge.\nFrom the user's standpoint, another shortcoming resides in the fact that both types of devices are inconvenient and time-consuming to mount to the wall and equally troublesome to move or remove. The problem is exacerbated by the recurrent need to move or replace the device to accommodate a different model or models.\nIn most instances, the existing wall-mounting devices for larger models provide only vertical support for the model. With the leading edge of each wing resting on an open bracket, the model is free to move from side to side, and to a lesser but significant degree, to rotate in pitch and move from front to rear within the bracket. The potential for structural damage is magnified by the size and weight of the model. In most instances, to modify the existing wall-mounts to limit or prevent such movement would require most of them to be custom-fit to each model. Given the problems mentioned in connection with the production and marketing of the basic devices themselves, the diversity of such improvements would make the manufacture, stocking, distribution, and inventorying prohibitively difficult and expensive.\nViewed against this background, an object of the subject invention is to provide a novel, unitary wall-mounted bracket for storing and displaying model airplanes vertically and flush with a wall, which minimizes both the wall area and the room space required for any given model, reduces the model's exposure to damage from external sources, and virtually eliminates the possibility of wing warp or internal damage to the fuselage due to gravity.\nAnother object is to provide a novel construction for a wall-mounted support that is sturdy, lightweight, and durable, yet can be installed and removed with a minimum of effort.\nYet another object is to provide a storage device for model airplanes which can be readily adjusted to conform to models with a variety of sizes and shapes.\nAn additional object is to provide an adjustable model airplane support that incorporates means for restricting the movement of the mounted model and thereby minimizing the likelihood of its accidental damage.\nStill another object is to provide a wall-mounted support device of the type described which can be fabricated from a minimum number of parts and modular components.\nA further object is to provide a wall-mounted model airplane storage device construction satisfying all of the foregoing objects that is comparatively inexpensive to manufacture and requires a minimum of maintenance and repair.\nYet a further object is to provide a wall-mounted support construction of this type that permits the standardization and use of interchangeable parts and modular components.\nA particular object of the invention is to provide a wall-mounted storage device for model airplanes, which incorporates in a single integrated structure adjustable means for supporting the weight of the model on the nose of the model.\nAnother particular object is to provide in such a device wing-retaining means for aiding in maintaining the model balanced upright on its nose.\nA further particular object is to provide in an alternative embodiment of the invention wing-retaining means which in addition to aiding in maintaining the model balanced upright support a portion of the weight of the model as well.\nStill another object is to provide a wall-mounted storage device for model airplanes, which incorporates in a single integrated structure adjustable means for supporting the weight of the model on both the wings and the nose of the model.\nYet a further object is to provide a wall-mounted support device of the type described in which the weight of the model is distributed fairly evenly between the wings and the nose of the model.\nA still further object is to provide a wall-mounted airplane model support device of this type having means for adjusting the horizontal and vertical positioning, the orientation, and the angular alignment of the wing-supporting means with respect to the nose-supporting means.\nOther objects will become apparent from the following summary of the invention and detailed description of its preferred embodiments."}
{"text": "1. Field of the Invention\nThe present invention relates to magnetic disk drives, and more particularly to a mechanism for preventing damaging contact between a magnetic head slider assembly and a rigid magnetic disk due to shock loads applied to a magnetic disk drive.\n2. History of the Prior Art\nAmong the better known data storage devices are magnetic disk drives of the type in which a magnetic head slider assembly floats on an air bearing at the surface of a rotating rigid magnetic disk. Such disk drive are often of the so-called Winchester type in which one or more rigid magnetic disks are located within a sealed chamber together with one or more magnetic head slider assemblies. The magnetic disk drive may include one, two or more rigid magnetic disks, and the slider assemblies may be positioned at one or both of the opposite sides of each of the magnetic disks.\nEach magnetic head slider assembly in magnetic disk drives of the type referred to is typically coupled to the free outer end of a different one of a plurality of elongated arms or load beams. The slider assembly is mounted in a manner to permit gimballed movement at a free outer end of the arm so that an air bearing between the slider assembly and the surface of the rigid magnetic disk can be established and maintained. The elongated arm is coupled to an appropriate mechanism for moving the arm across the surface of the disk so that a magnetic head contained within the slider assembly can address different ones of concentric data tracks on the disk for writing information into or reading information from the data tracks.\nAn example of an arm assembly having a gimballed mount for a magnetic head slider assembly is provided by U.S. Pat. No. 3,931,641 of Watrous. The arm assembly described in the Watrous patent includes a relatively rigid load beam having a rigid bearing member at a free outer end thereof for receiving a protuberance on a spring element. The spring element is spot welded to the load beam and has an end thereof defining a flexure. The flexure includes a pair of stiff crosslegs mounted on an opposite pair of flexible outer fingers and a central finger. The central finger mounts a magnetic head slider assembly, and gimballed movement is provided by the load protuberance on the spring element which is held in contact with the bearing member at the end of the rigid load beam. Such arrangement provides desired gimballing action by allowing a pitch and roll of the slider assembly around mutually orthogonal axes while at the same time resisting radial, circumferential and yaw motions.\nIn magnetic disk drives of the type described, physical contact of the surface of the magnetic disk containing the data tracks by the slider assembly must be avoided. Consequently arrangements must be provided for preventing such contact when the magnetic disk is at rest or is otherwise not rotating at its nominal operating speed. When the disk is rotating at its nominal operating speed, the air bearing between the slider assembly and the disk is usually sufficient to prevent contact therebetween. One technique commonly employed is to move the slider assembly onto a portion of the disk where physical contact can be tolerated whenever the disk is decelerated to rest. One or more so-called parking zones are provided on the surface of the disk where the slider assembly may rest. When the disk is accelerated to its nominal operating speed, the air bearing again forms and the slider assembly may then be moved to the data tracks. Another common technique for preventing contact between the slider assembly and the disk surface when the disk is decelerated to rest is to pivot the arm so as to move the slider assembly away from the surface of the disk. This moves the slider assembly into an unloaded position well away from the disk surface. When the disk is again accelerated to the nominal operating speed, the arm is pivoted so as to bring the slider assembly back into a loaded position in which the air bearing forms between the slider assembly and the surface of the disk.\nAn example of an arrangement for raising and lowering the arm so as to move the slider assembly between loaded and unloaded positions is provided by a copending application of Warrent L. Dalziel, Ser. No. 759,900, filed July 29, 1985 and commonly assigned with the present application. The Dalziel application describes a mechanism for loading and unloading one or more slider assemblies using simple mechanical apparatus. Such apparatus engages the arms to raise them into unloaded positions and at the same time limit their radial movement. A separator element is movable between first and second positions to selectively engage the arms. When the separator element is in a first position, the slider assemblies are in loaded positions, and radial motion of the slider assemblies is not inhibited. Movement of the separator element into a second position moves the slider assemblies into unloaded positions and at the same time locks an associated carriage assembly in a selected position to inhibit radial motion of the slider assemblies.\nEach time the disks are decelerated to rest in the arrangement described in Dalziel application Ser. No. 759,900, the slider assemblies are raised into the unloaded positions and at the same time radial movement of the arms is limited. This does much to protect the sensitive components of the magnetic disk drive against damage due to shock loads such as may occur if the disk drive is dropped or otherwise bumped. Nevertheless, depending upon the manner of coupling the slider assemblies to the ends of the elongated arms, it may be difficult or impossible to prevent damaging contact between the slider assemblies and the disk surfaces. This is so even though the arms are pivoted to move the slider assemblies into unloaded positions well away from the disk surfaces. Inadvertent dropping of the magnetic disk driven onto a hard surface or other application of severe shock loads may still result in the slider assemblies contacting the disk surfaces with possible damage to the slider assemblies or the disk surfaces or both. Typically the slider assembly is coupled to the free outer end of the arm by a flexure element which is highly flexible and which therefore permits substantial movement of the slider assembly away from the arm in response to shock loads and other substantial forces.\nAccordingly, it is an object of the invention to provide an improved magnetic disk drive, and particularly an improved disk drive of the type in which one or more magnetic head slider assemblies are employed in conjunction with one or more rigid magnetic disks.\nIt is a further and more specific object of the invention to provide an improved magnetic disk drive in which the possibility of damaging contact of the disk surfaces by the slider assemblies in response to shock loads when the arm is raised into an unloaded position is greatly minimized or eliminated."}
{"text": "I. Field of Invention\nThe invention generally relates to wireless communication and, more particularly to an orthogonal frequency division multiplexing scheme to modulate digital data into audible tones for wireless communication and authentication.\nII. Description of the Related Art\nAdvances in communication technology have made it easier and faster to share and/or transfer information. High volumes of data can be communicated through data transmission systems such as a local or wide area network (for example, the Internet), a terrestrial communication system or a satellite communication system. These systems require complicated hardware and/or software and are designed for high data rates and/or long transmission ranges.\nFor transfers of data over short distances, such as between a personal computer and personal data assistants (PDAs), the systems above may be an inconvenient communication medium for users because of the complexity, delay and often the cost involved in accessing the systems. Accordingly, various communication systems have been developed using communication medium such as radio frequency (RF) or Infrared (IR) to transmit data. However, these systems also require specialized communication hardware, which can often be expensive and/or impractical to implement. Non-wireless connections can also be used to transfer data. However, to use non-wireless types of connections, users must physically have as well as carry wires or cables and make the physical connections for communication. This can be burdensome and inconvenient to users.\nMoreover, with the increase in electronic commerce, the opportunities for fraudulent activity have also increased. Misappropriated identity in the hands of wrongdoers may cause damage to innocent individuals. In the worst case scenarios, a wrongdoer may actually purloin a party's identity in order to exploit the creditworthiness and financial accounts of an individual.\nIn order to prevent unauthorized persons from intercepting private information, various security and encryption schemes have been developed so that private information transmitted between parties is concealed. However, concealment of private information is only one aspect of the security needed to achieve a high level of consumer confidence in electronic commerce transactions. Another aspect is authentication.\nAccordingly, electronic authentication of an individual may currently be performed by authentication through knowledge, such as a password, secure identification or a personal identification number (PIN); authentication through portable objects, such as a credit card, or a proximity card; and/or authentication through personal characteristics (biometrics), such as fingerprint, DNA, or a signature.\nWith current reliance on electronic security measures, it is not uncommon for an individual to carry multiple authentication objects or be forced to remember multiple passwords. Authentication through knowledge can thus be problematic for individuals who are forced to remember multiple passwords or PINs. Accordingly, such information may be written down on objects. However, writing down such information leaves an individual vulnerable to the theft of passwords, PIN codes, or even the object itself. Accordingly, there may still be a burden and/or expense to conceal the object from others.\nTherefore, there is need for a less complex, yet user-friendly, inexpensive and/or efficient way to share and/or transfer information."}
{"text": "1. Field of the Invention\nThe present invention relates to a water-cooled type internal combustion engine mounted to a farming motor vehicle such as a tractor and a combine harvester and more particularly to a device adapted to cool a muffler of such an engine by means of a radiator fan.\n2. Description of the Prior Art\nConventionally an engine for a farming motor vehicle is provided with a muffler which is projected upward largely from an engine room of the motor vehicle. In recent years in an attempt to improve an external appearance of the motor vehicle, the muffler is mounted within the engine room.\nAs such a muffler of the type mounted within the engine room has been known the one disclosed in Japanese Provisional Utility Model Publication No. 1986-32510.\nHere, in the engine room extending in the fore and rear direction of the motor vehicle there are arranged a radiator, a radiator fan and an engine body in order from foreside while a muffler body is arranged in a space on the lateral side of the engine body so as to extend in the fore and rear direction. An exhaust outlet pipe is projected downward from the upwind side end portion of the muffler body, and a muffler cover is mounted to the muffler body with a cooling air passage provided therebetween for cooling the muffler. A portion of cooling air flow delivered by means of the radiator fan is introduced into the muffler cooling air passage so as to pass therethrough.\nIn the aforementioned conventional embodiment there is the advantage in that the muffler body is adapted to be cooled by means of the cooling air flow delivered from the radiator fan so as to prevent the temperature within the engine room from being raised by a heat of the muffler body and to cool down the temperature of an exhaust gas discharged from the outlet pipe.\nThere is, however, a following disadvantage associated with the conventional embodiment. That is, since it is required to increase a quantity of the cooling air flow delivered by means of the radiator fan when the temperature of the exhaust gas from the outlet pipe is intended to be further cooled down, the radiator fan becomes large in size and resultantly the engine itself becomes large in overall size."}
{"text": "The present invention relates to new reactive systems comprising organic polyisocyanates, organic polyepoxides, polymer modifiers, special catalysts and optionally other auxiliary agents and additives.\nReactive systems based on epoxide/isocyanate (EPIC) components are described, for example, in DE 36 44 382.\nEPIC resins have good processing and curing properties and when cured give rise to articles having good mechanical strength, chemical resistance and high glass transition temperatures. However, the products are relatively brittle. The same applies to epoxide resins and duroplasts, in general.\nIt has long been the wish of users of such products to have available duroplastic synthetic resins which combine high rigidity with chemical resistance, heat resistance, improved toughness and ductility.\nIt was, therefore, an objective of this invention to develop EPIC resins which when cured give rise to articles with improved toughness without loss of rigidity and heat resistance. The present invention provided EPIC resins which meet this objective by employing polymer modifiers. It is known that modifiers are used as polymer blending partners for thermoplastic synthetic resin blends (e.g. on the basis of polyvinyl chloride (PVC). It was, however, surprising that polymer modifiers of the type of graft rubbers achieve the objective of this invention."}
{"text": "Data reading devices, such as barcode or optical code scanners, RFID readers, and the like, are a popular means for data acquisition in computerized processing systems. Barcode scanners are used to optically capture barcode patterns or other symbols or information imprinted on different surfaces in order to transmit the information encoded in the barcode pattern or symbol to a host processing device.\nTwo major types of commonly used barcode scanners are flying spot laser scanners and image based barcode scanners. Flying spot laser barcode scanners generally obtain barcode information by sweeping a laser spot across the barcode. The laser spot may be generated from a light source inside an oscillating reflecting surface, typically a mirror. The light reflected from the barcode is collected by a photosensor, which outputs an analog waveform representing the relative spacing of the bars in the barcode. The analog signal may then be digitized and decoded into data representing the information encoded in the barcode.\nBarcode scanning devices based on solid state image circuitry, such as charge coupled devices (CCDs), are also conventionally known. These types of barcode scanners are typically implemented using either a one-dimensional or two-dimensional imaging array of photosensors (or pixels) to capture the barcode. One-dimensional CCD readers capture a linear cross section of the barcode at once, producing an analog waveform whose amplitude represents the relative darkness and lightness of the barcode. Two-dimensional CCD readers capture an entire two-dimensional image at once.\nA relatively new type of data reading technology is radio frequency identification (RFID). An RFID system typically employs at least two components, a “transponder” or “tag,” which is attached to the physical item to be identified, and a “reader,” which sends an electromagnetic signal to the transponder and then detects a response. Typically, the reader emits an RF signal which is received by the transponder after the transponder comes within an appropriate range. In response to the signal from the reader, the transponder sends a modulated RF signal back to the reader. The reader detects this modulated signal and receives the information from the transponder by decoding the modulated signal. After receiving the information from the transponder, the reader can either store the decoded information or transmit the decoded signal to a computer or other host.\nThe transponder used in an RFID system may be either “passive” or “active.” A passive transponder can be a simple resonant circuit, including an inductive coil and a capacitor. Passive transponders are generally powered by the carrier signal transmitted from the reader. Active transponders, on the other hand, generally include transistors or other active circuitry, and require their own power source.\nBarcode scanners and RFID readers may be physically separate from host processing devices, such as personal computers, which provide decoding software, computer storage, and communication interfaces. Increasingly, however, barcode scanners and RFID scanners are being integrated into portable host processing devices. For instance, the Falcon® 4420 data collection terminal available from Datalogic Mobile, Inc. of Eugene, Oreg., combines a laser barcode scanner, storage memory, and wireless communication interface in a compact, portable form factor. Such portable devices can be easily transported to different locations within a warehouse, retail store, or the like, greatly simplifying the tasks of receiving and inventorying.\nUpdating configuration settings, updating an operating system (OS) image, and booting an alternate OS on a portable data reader may require intervention by the user (e.g., to initiate updates, monitor status, and respond to prompts) and may also require connection to a host computer. Accordingly, the present inventor has recognized a need for techniques to update configuration settings, update an OS image, and boot an alternate OS on a portable data reader that does not require a host computer and that requires little or no user intervention."}
{"text": "During operation a fault, such as a fuel leak, may occur in a power plant.\nIt would be beneficial to be able to detect faults such as fuel leaks during operation of a power plant so that an operator of the power plant may be alerted to the fault."}
{"text": "Materials which abrade readily in a controlled fashion are used in a number of applications, including as abradable seals. Very few thermal spray abradable coatings, however, are suitable for high-temperature applications. In general, contact between a rotating part and a fixed abradable seal causes the abradable material to erode in a configuration which closely mates with and conforms to the moving part at the region of contact. In other words, the moving part wears away a portion of the abradable seal so that the seal takes on a geometry which precisely fits the moving part, i.e., a close clearance gap. This effectively forms a seal having extremely close tolerances.\nOne particular application for abradable seals in high-temperature environments is their use in axial flow gas turbines. The rotating compressor or rotor of an axial flow gas turbine consists of a plurality of blades attached to a shaft which is mounted in a shroud. In operation, the shaft and blades rotate inside the shroud. The inner surface of the turbine shroud is most preferably coated with an abradable material. The initial placement of the shaft and blade assembly in the shroud is such that the blade tips are as close as possible to the abradable coating.\nAs will be appreciated by those skilled in the art, it is important to reduce back flow in axial flow gas turbines to maximize turbine efficiency. This is achieved by minimizing the clearance between the blade tips and the inner wall of the shroud. As the turbine blades rotate, however, they expand somewhat due to the heat which is generated. The tips of the rotating blades then contact the abradable material and carve precisely defined grooves in the coating without contacting the shroud itself. It will be understood that these grooves provide the exact clearance necessary to permit the blades to rotate at elevated temperatures and thus provide an essentially custom-fitted seal for the turbine.\nIn other gas turbines, the initial clearance is somewhat greater and the abradable coating is intended to protect the shroud and blade tips against wear during transient conditions (e.g., power surges).\nIn order for the turbine blades to cut grooves in the abradable coating, the material from which the coating is formed must abrade relatively easily without wearing down the blade tips. This requires a careful balance of materials in the coatings. In this environment, an abradable coating must also exhibit good resistance against particle erosion and other degradation at elevated temperatures. As known by those skilled in the art, however, few conventional thermal spray abradable coatings have the desired high-temperature performance characteristics.\nA number of abradable coatings are known in the art. Limited success has been achieved by others with the use of ZrO.sub.2 based ceramic coatings in abradable applications. ZrO.sub.2 based powders have also been blended with plastic based powders, the blended mixture being plasma sprayed to form abradable coatings. These approaches, however, have produced coatings which exhibit limited abradability at high temperatures. In addition, the plastic powders tend to degrade during thermal spraying, producing inconsistent microstructures and inferior abradability.\nOther conventional abradable coatings include such cellular or porous metallic structures as those illustrated in U.S. Pat. Nos. 3,689,971, 4,063,742, 4,526,509, 4,652,209, 4,664,973, and 4,671,735. Low melting point metallic coatings of indium, tin, cadmium, lead, zinc, and aluminum alloys have been suggested for use in providing \"ablative\" seals wherein heat generated by friction melts a clearance gap in the coating. This approached is exemplified in U.S. Pat. Nos. 2,742,224 and 3,836,156. Ceramics such as ZrO.sub.2 and MgO for use in forming abradable coatings are also shown in U.S. Pat. Nos. 4,405,284, 4,460,311, and 4,669,955.\nIn U.S. Pat. Nos. 3,508,955, a composite material is disclosed which comprises a porous metal impregnated with a fluoride of metals selected from Groups I and II of the Periodic Table of the Elements. The use of fluoride salts and a barium fluoride-calcium fluoride eutectic is specifically mentioned as is the use of the material in bearings and seals. It is also disclosed therein that the resultant material can be sprayed with a surface layer of fluoride eutectic slurry which is then dried and sintered.\nIn U.S. Pat. No. 4,867,639, abradable coatings for use in turbine or compressor shrouds are disclosed which are described as low melting fluoride compounds such as BaF.sub.2, CaF.sub.2 and MgF.sub.2 incorporated into a higher melting temperature ceramic or metallic matrix. It is disclosed that, alternatively, the soft ceramic phase may be used to fill or impregnate a honeycomb shroud lining made of the higher melting temperature ceramic or metal alloy, so that the soft ceramic is not eroded by hot gases in the turbine. Zirconia and/or alumina are disclosed as the preferred high melting temperature ceramic, and NiCr and NiCrAl are disclosed as preferred metals.\nThe use of metal matrix coatings having a plastic component such as a polyimide are also known for use in forming an abradable seal in high-efficiency compressors. Due to the lower temperatures generated in the compressor and the fact that the rotating blades are generally softer than those found in the turbine section, plastics have been used in lieu of solid lubricants such as CaF.sub.2. While the lower melting point of plastics is advantageous in such low temperature applications, the use of these coatings has not been successful in high temperature applications.\nIn U.S. Pat. No. 5,196,471, \"Thermal Spray Powders for Abradable Coatings Containing Solid Lubricants and Methods of Fabricating Abradable Coatings,\" thermal spray powders are described which are characterized by the presence of a matrix-forming component, a solid lubricant component and a plastic component. Abradable coatings formed by thermal spraying the powders abrade readily to form abradable seals. The abradable coatings have a metal, metal alloy, or ceramic matrix with discrete inclusions of solid lubricant and plastic. Therein, the use of Zirconia is described as a preferred ceramic for use as the matrix-forming component.\nTherefore, it would be desirable to provide a composite material which abrades readily at high temperatures without producing significant wear of rotating parts.\nIt would also be desirable to provide such a material which can be fabricated using conventional thermal spray techniques.\nIt would still further be desirable to provide a coating for forming abradable seals which can be custom formulated for a particular operating environment.\nThe present invention achieves these goals by providing thermal spray powders which are a two component blended mixture that forms high-temperature, abradable coatings by conventional thermal spray application."}
{"text": "Enterprise content management (ECM) covers a broad range of applications, including document management (DM), Web content management (WCM), records management (RM), digital asset management (DAM), search of managed content, and the like. A content management system (CMS) suitable for managing the various content (also referred to herein in some examples as “files” or “documents”) that an enterprise produces or generates, retains or otherwise stores, manipulates or modifies, etc. can support the requirements of one or more of such applications, and optionally other requirements, to provide a coherent solution in which content processes, management processes, and the like are capable of accessing content across a variety of applications subject to access controls, permissions, and the like. Content managed by a CMS can include one or more of documents, images, photos, Web pages, records, XML documents, other unstructured or semi-structured files, etc. Content retained in a CMS can also include directory structures such as folders, file trees, file plans, or the like, which can provide organization for multiple content items in addition to storing or otherwise representing relationships between content item, etc. An “enterprise” can generally refer to an organization, such as for example a business or company, a foundation, a university, or the like, and can have content requirements related to one or more business processes, content uses, etc.\nA CMS manages the actual digital binary content, the metadata that describes a context of each content item, associations between a content item and other content or content items, a place and classification of a content item in a repository, indexes for finding and accessing content items, etc. The CMS can also manage processes and lifecycles of content items to ensure that this information is correct. The CMS can also manage one or more workflows for capturing, storing, and distributing content, as well as the lifecycle for how long content will be retained and what happens after that retention period.\nA CMS for use in enterprise content management can include one or more of document management tools, applications, and interfaces to support general office work, search, and discovery. Workflow management capabilities of a CMS can support numerous business processes, optionally including, but not limited to, case management and review and approval. Collaboration applications and services of a CMS can support the collaborative development of information and knowledge in the creation and refinement of content and documents. This collaborative development of information and knowledge can be achieved through providing access to content managed by the CMS to multiple users. To prevent conflicting or discontinuous editing streams, a user can be allowed to check out or lock content for modification and check in the modified content such that other users are prevented from editing content concurrently. Web content management services of a CMS, which can be scalable, can support the delivery and deployment of content from the enterprise to its customers. Records management capabilities of a CMS can capture and preserve records based upon government-approved or other standards. A standards-based platform can also provide access to applications that use these standards, such as publishing, image management, email management, etc."}
{"text": "The term computer “hacker” is increasingly used to refer to persons who engage in malicious or illegal activities to gain access to, or attack computer systems without authorization. Such activities by hackers have evolved far beyond those that were simply pranks or caused some minor inconveniences into a key component of highly organized criminal enterprises in which billions of dollars can be made each year.\nHackers often seek to launch attacks on computer systems in an automated manner by using large networks called “botnets” of compromised computers called “bots” (i.e., software robots) or “drones.” While bots can be supported by hosts that bypass most local Internet content regulation (so called “bullet-proof hosts”), bots are primarily found in computers used by innocent home users who are completely unaware that their systems have been taken over and are being used for illegitimate purposes. Botnets are thought to be organized in herds as large as one hundred thousand to a half million or more bots that can be geographically spread over many countries.\nBotnets can employ both active and passive attacks. In an active attack, a botnet originates attacking traffic such as spam, adware, or denial of service (“DoS”) traffic which is sent over a network such as the Internet to its victims. In a passive attack, bots function as servers which, when accessed by a user, serve malware such as viruses, rootkits, trojan horses etc., typically using HTTP (Hypertext Transfer Protocol).\nReputation services have been established to address the problem of automated attacks and other hacker activities by compiling black lists of URLs (Uniform Resource Locators) and IP (Internet Protocol) addresses of known adversaries. A variety of technologies such as mail relay servers, firewalls, and unified threat management gateways can query the reputation service through an online connection to decide whether to accept traffic from, or send traffic to, a given computer on the Internet.\nCurrent reputation services often run their own laboratories that are equipped with a variety of tools which are used to scan the Internet to locate adversaries and establish the reputation. These tools include web crawlers, honeypots (passive, dummy data or network sites that appear to contain information of value to attract attackers), honey monkeys (virtual computers that visit websites and seek code designed to attack a computer), virtual machines, and other global sensors.\nReputation services face several significant challenges that can affect their use and success in combating hackers. For example, reputation services must reliably detect and confirm adversaries that are deployed in vast numbers all over the world. Hackers can also change URLs and IP addresses of bots frequently, so reputation services must be able to dynamically respond with equal speed to detect them and not block legitimate users who might reuse the same URL or IP address a few hours later. This problem of false positives in which URLs and IP addresses of innocent (i.e., non-malicious) computers are wrongly identified as adversaries can cause significant disruptions to users and result in high costs to service providers to resolve disputes and restore services.\nIn addition, reputation services need to successfully deal with hackers who are increasingly targeting attacks on small sets of customer or enterprise networks that frequently go unobserved by the technologies of existing reputation services.\nThis Background is provided to introduce a brief context for the Summary and Detailed Description that follow. This Background is not intended to be an aid in determining the scope of the claimed subject matter nor be viewed as limiting the claimed subject matter to implementations that solve any or all of the disadvantages or problems presented above."}
{"text": "The present invention relates to a semiconductor integrated circuit and an IC card, and more particularly to a technique that can be effectively applied to, for instance, the prevention of reverse engineering of a cryptographic key or the like held by a semiconductor integrated circuit, such as a microcomputer for IC cards.\nAlong with the development of semiconductor technology, it has become a common practice to settle accounts in a safe and secure way by incorporating integrated circuits (ICs) into credit cards, securities or the like and communicating account information in an encrypted form. This IC-based method makes it more difficult to forge a paper or use another person's name than the conventional method of using magnetic records, and as such is beneficial to both end users and service providers.\nCryptographic algorithms have been studied for many years, making it extremely difficult to infer a cryptographic key from signals obtained by tapping a communication line, and this risk is virtually negligible. A real problem, however, is posed by attempts to open an IC and reverse-engineering it to directly read internal information or a cryptographic key contained in the IC. Reverse engineering is a technique or an act to analyze the structure and/or specifications of a hardware or software product and thereby and thereby reveal the technical information contained.\nPreviously devised reverse engineering techniques was to read internal information or a cryptographic key contained in an IC card by supplying a clock of an illegitimate frequency to the card, abruptly raising or reducing the voltage of power supply to it or irradiating it with a powerful electromagnetic wave to abnormally operate the IC card. On the part of the IC card, such intrusions were fought back by detecting such abnormal states, and preventing internal information or the cryptographic key from being read out on the basis of the detected acts.\nFor instance, according to a technique described in Patent Reference 1, an IC chip for an IC card is provided with an unsealing sensor and, if it detects unsealing, a CPU will erase data in the memory to ensure safe protection of secrets.\nPatent Reference 2 describes a technique according to which a small hole is bored in a part of a package to seal and shadow the circuit configuration so that only the sensor part for light detection be illuminated with a light and the circuit can operate normally only when the light is detected. Since an unauthorized analyst would open the package in a dark place to avoid adverse impacts of light, the circuit would operate differently from its usual way in a state in which no light is detected. This different operation makes impossible analysis of the normal operation of the circuit and accordingly illegitimate reading of the stored information.\nPatent Reference 3 discloses a technique according to which a plurality of light receiving elements are integrated on an IC in a distributive way, and each of the plurality of light receiving elements is connected to one of a connection line connected to a nonvolatile memory cell, a connection line connected to a logic circuit and a connection line connected to a logic element and, by intercepting, establishing conduction of or grounding these connection lines, the circuits related to the respective connection lines are inhibited from normally operating so that the internal information contained in the IC can be protected even if it is unsealed.\nPatent Reference 1: Japanese Published Unexamined Patent Application No. Hei 10(1998)-320293\nPatent Reference 2: Japanese Published Unexamined Patent Application No. 2000-216345 (paragraphs 0009 through 0011)\nPatent Reference 3: Japanese Published Unexamined Patent Application No. Hei 11(1999)-102324"}
{"text": "1. Technical Field\nThe present invention relates to a wireless communication device.\n2. Background Technology\nAs a method for establishing a network, there is a method for performing Wi-Fi Direct (hereinafter referred to as “Wi-Fi Direct Connection”) under Wi-Fi Alliance standard development. A wireless communication device that is possible to perform the Wi-Fi Direct Connection can directly communicate with other wireless communication devices that are also possible to perform the Wi-Fi Direct Connection as one-on-one communication or a multiple devices communication without using an access point. Also, it is possible to communicate with a wireless communication device that is possible to perform a wireless LAN connection method used with Wi-Fi (hereinafter referred to as “Wi-Fi Connection”) function as an access point. In addition, since the Wi-Fi Direct Connection uses the Wi-Fi technologies, the high-speed communication is possible. By using the Wi-Fi Direct Connection, for example, it is possible to directly send an image to a printer from a camera or a cellular phone and print it out.\nJapanese Laid-open Patent Publication No. 2011-166417 (Patent Document 1) and http://205.149.128.22/Wi-Fi_Direct.php (Non-patent Document 1) are examples of the related art."}
{"text": "1. Field of the Invention\nThis invention relates to a new and novel process for the synthesis of compounds heretofore known only as naturally occurring substances. Specifically, these compounds were heretofore isolated from Steganotaenia araliacea, a plant native to Ethiopia, by extraction with alcohol.\n2. Description of the Prior Art\nA. Steganacin and Steganangin are described as possessing anti-leukemic activity in the article of Kupchan, Britton, Ziegler, Gilmore, Restino and Bryan, Journal of the American Chemical Society, 95 (4), 1335 (1973). The article also describes the compounds known as Steganone and Episteganol."}
{"text": "The present invention relates to bearing control and more specifically to bearing wear detection.\nIn many applications, bearing failure can lead to costly damage to the machines in which the bearings are operating and provides a potential for breakdown at unknown and often inconvenient intervals. Therefore, devices for monitoring bearing wear have shown themselves to be indispensable. Bearings for moving machinery parts and shafts are usually securely covered and sealed for reasons of safety, to prevent foreign material from contaminating the bearing surfaces, and to contain lubricating fluids used in operation of the bearings. In general, shaft bearings are protected by oil or watertight seals and in some cases these seals cannot be removed while the shafts are still in operation or the seals are under high pressure fluid. Under such conditions, if the bearing clearances have to be measured to assure the operation of the machinery, the operation of the shaft has to be shut down, and the high fluid pressure behind the seals, if any, has to be brought down to a workable level before the bearing clearance measurements can be performed. It is known that bearing clearances of submarine control surfaces are measured by means of hydraulic jacks, lift from cranes, and dial gauges. However there were cases, utilizing the prior art measuring procedures, in which the bearings were opened up unnecessarily for repair since the bearings were subsequently found to be in good condition. To obviate the need for such time consuming and involved procedures various systems have been devised to remotely ascertain bearing clearances.\nPrior art apparatus used to remotely monitor bearing wear operate in accordance with the principle of galvanic currents passing through the shaft and bearing whereby the illumination of a small lamp on metallic contact taking place between bearing and shaft serves as an indication of the critical condition. Variations on this system include the use of sets of wires annularly disposed around the bearing shaft so that any eccentricities in the rotation of the shaft induced by bearing wear lead to electrical contact which in turn triggers an alarm indicative of excessive bearing wear. However, in all of these systems, the input of electric current to the shaft requries contact points with the latter, e.g. brushes made of speical copper-graphite in contact with a needle mounted centrally on the shaft. This device has shown itself to be inadequate in view of the unavoidable material wear and in the inconsistancy of contact resistances.\nAn attempt has also been made to determine the magnitude of change in the lubrication gap between shaft and bearing with capacitive measuring devices. This system can only be used in special cases as a contact indicator, namely only in the case where the measuring device is situated at the same height as the load surface of the bearing and the if shaft touches the bearing at this point. Further, arrangements have been proposed whereby capacitive measurements or the clearance of a rotating body take place with the aid of fixed electrodes. With this arrangement, two electrodes are screened off from one another and an alternating voltage of constant amplitude and frequency is fed to one of the electrodes and the voltage arising at the second electrode due to capacitive coupling regulated by the rotating body is amplified and led to an indicating instrument. This system generally works well except that it does not allow the position or degree of the bearing wear to be remotely determined, it requires expensive equipment, and is subject to malfunction if impurities exist in the lubricating fluid."}
{"text": "Bednorz and Muller, Z. Phys. B64, 189 (1986), disclose a superconducting phase in the La-Ba-Cu-O system with a superconducting transition temperature of about 35 K. The presence of this phase was subsequently confirmed by a number of investigators [see, for example, Rao and Ganguly, Current Science, 56, 47 (1987), Chu et al., Science, 235, 567 (1987)]. Chu et al., Phys. Rev. Lett., 58, 405 (1987), Cava et al., Phys. Rev. Lett., 58, 408 (1987), Bednorz et al., Europhys. Lett., 3, 379 (1987)]. The superconducting phase has been identified as the composition La.sub.l-x (Ba,Sr,Ca).sub.x O.sub.4-y with the tetragonal K.sub.2 NiF.sub.4 -type structure and with x typically about 0.15 and y indicating oxygen vacancies.\nWu et al., Phys. Rev. Lett., 58, 908 (1987), disclose a superconducting phase in the Y-Ba-Cu-O system with a superconducting transition temperature of about 90 K. The compounds investigated were prepared with nominal compositions (Y.sub.1-x Ba.sub.x).sub.2 CuO.sub.4-y and x=0.4 by a solid-state reaction of appropriate amounts of Y.sub.2 O.sub.3, BaCO.sub.3 and CuO in a manner similar to that described in Chu et al., Phys. Rev. Lett., 58, 405 (1987). This reaction method comprised heating the oxides in a reduced oxygen atmosphere of 2.times.10.sup.-5 bars (2 Pa) at 900.degree. C. for 6 hours. The reacted mixture was pulverized and the heating step was repeated. The thoroughly reacted mixture was then pressed into 3/16 inch (0.5 cm) diameter cylinders for final sintering at 925.degree. C. for 24 hours in the same reduced oxygen atmosphere.\nHundreds of other papers have since disclosed similar solid state reaction processes. Other papers have disclosed various solution and precipitation methods for preparing the reactants to be heated at temperatures of 800.degree.-850.degree. C. and above.\nHirano et al., Chemistry Letters, 665, (1988), disclose a process for producing Y-Ba-Cu-O superconductors by the partial hydrolysis of a solution of Ba metal, Y(O-iPr).sub.3 and Cu-acetylacetonate or Cu-alkoxides in 2-methoxy or 2-ethoxy ethanol. The solution was stirred in dry nitrogen and heated at 60.degree. C. for 12 hours. The solution was then hydrolyzed by the slow addition of water diluted with solvent. Stirring and heating continued for several hours. Stirring continued while the solution was evaporated under vacuum at about 60.degree. C. and an amorphous precursor powder was obtained. The powder was calcined in flowing oxygen at temperatures between 800.degree. and 950.degree. C. for up to 24 hours. The calcined powder was pressed and sintered in flowing oxygen at temperatures up to 920.degree. C. and then annealed at temperatures between 450.degree. and 550.degree. C.\nIt is highly desirable to form precursors that can be used to produce powders that have small size particles, i.e., generally sub-micron in size, and that can be pressed into desired shapes, sintered and converted to superconducting Y-Ba-Cu-O."}
{"text": "1. Field of the Invention\nThis invention relates to the use and display of a drill bit, probe, or other elongate instrument with a computer assisted surgery system.\n2. Description of the Related Art\nComputer assisted surgery systems have been developed to aid a surgeon in more accurately positioning instruments during a surgical procedure. These systems have been described for use with CT images, fluoroscopic images, and images derived from other modalities as well as without images or in conjunction with a surgical robot. When using images of a patient, they are often referred to as image guided surgery systems. These systems typically use a localizing device to track surgical instruments in real time. One commonly used localizing device is an optical localizer that employs a stereoscopic camera system to view infrared light emitters or reflectors that are placed on the surgical instruments. Other localizing devices may use electromagnetic radiation or a passive manipulator arm with position encoders at the joints. The computer assisted surgery system then uses the information regarding the instrument\"\"s position to superimpose a representation of the instrument over the images or to generate some other graphic or alphanumeric display.\nCommonly, surgical procedures performed with computer assisted surgery systems involve the insertion into the patient of an elongate instrument such as a drill bit or a probe. In these cases the computer-asssistance may consist of superimposing a line representing the trajectory of the drill or probe over images of the patient. Many systems additionally provide for a means for determining the location of the tip of the probe or drill bit. By comparing this tip location to the location of a target feature in the body, the software can calculate the distance remaining in the trajectory. This information can be displayed as a bar graph as taught by U.S. Pat. Nos. 5,251,127 and 6,006,127 or as a numeric value or other graphic as taught by U.S. Pat. No. 5,638,819. However, these approaches are not suitable for orthopaedic procedures, where the surgeon does not typically record a target point prior to the procedure.\nAnother problem with the approach of tracking the drill itself is that the drill bit may bend. This can cause the trajectory displayed by the computer assisted surgery system to be significantly different than the actual path of the drill bit. Bending of the drill bit is most likely to occur at a location between the drill and the body part, either as a result of deflection of the drill bit tip upon entering the body part or because of lateral forces placed on the drill by the surgeon. One way to alleviate this problem is to use a tool guide that allows the trajectory of the drill bit or probe to be measured as it enters the body part. Such tool guides are common in orthopaedic surgery and neurosurgery. Computer assisted versions of tool guides are described in U.S. Pat. Nos. 5,517,900 and 5,904,691, among others. With a tool guide placed flush against the bone the opportunity for drill bit bending is minimized and the trajectory displayed by the computer more closely matches the actual trajectory. However, tracking just the tool guide provides information regarding only the trajectory of the drill bit and not its progress into the body.\nThe present invention provides the surgeon with improved information for more accurately inserting a drill bit or other elongate object into the body, when used with a drill guide and a computer assisted surgery system.\nAccordingly, one objective of the present invention is to provide a computer assisted surgery system for accurately positioning a drill bit into a body part.\nAnother objective of the invention is to provide a technique and apparatus for calculating and graphically displaying the pose of a drill bit attached to a drill and passed through a drill guide by determining the poses of the drill and drill guide. The invention has the advantage of accurately displaying the orientation and location of the drill bit.\nStill another objective of the invention is to provide a technique and apparatus for reporting to the surgeon misalignment, malposition, or excessive bending of the drill bit relative to the drill guide.\nYet another objective of the invention is to provide a technique and apparatus for accurately calculating and graphically displaying the pose of a drill bit attached to a drill with only a single visible localizing emitter and passed through a drill guide.\nStill another object of the invention is to provide a surgical drill, with three or more localizing emitters integrated into the drill housing, whose pose can readily be determined with the use of a localizer.\nThese and other objects of the present invention are achieved by the use of a computer assisted surgery system, including a drill and a drill guide, each outfitted with three or more localizing emitters, and an optical localizer for measuring the pose of the drill and drill guide. A drill bit is attached to the drill and the location of the tip of the drill determined during a calibration step in which the tip of the drill is placed against a dimple in the drill guide and the poses of the two instruments recorded. During drilling, the drill bit is inserted through the bore of the drill guide and the poses of the drill guide and drill are recorded in real time. The pose of the drill bit is calculated from the orientation of the drill guide and the location of the drill. A graphic representation of the drill bit is then displayed on the monitor screen in an appropriate position relative to the stored images or other instruments.\nIn an alternative embodiment of the drill, the three or more localizing emitters are integrated into the drill housing and are spaced widely apart to facilitate accurate pose determination.\nIn an alternative embodiment of the computer assisted surgery system, the drill may be outfitted with one or more localizing emitters. In the calibration step, the distance between the drill\"\"s localizing emitter and the tip of the drill is determined. During drilling, the pose of the drill bit is calculated from the orientation of the drill guide and the location of the drill\"\"s localizing emitters, prior to display of its representation on the monitor screen."}
{"text": "Improvements in current practices of viral marker screening and donor self-exclusion are continuously increasing the safety of the blood supply. However, despite these practices, a risk of transmission of pathogens with the transfusion of cellular components of blood remains since current screening tests do not screen for rarely occurring or as yet unknown transfusion transmissible pathogens (Dodd, R. Y. New Engl. J. Med. 327:419-421 (1992); Soland, E. M., et al. J. Am. Med. Assoc. 274:1368-1373 (1995); Schreiber, G. B., et al. New Engl. J. Med. 334:1685-1690 (1996)).\nTo combat the deficiencies associated with screening practices, the use of sterilization procedures of blood, red blood cell concentrates (RBCC), and other blood-derived components hold promise for eliminating pathogen transmission. In this connection, various approaches have been used to sterilize blood cells, the most efficacious so far use photochemical methods (Ben-Hur, E. and B. Horowitz Photochem. Photobiol. 62:383-388 (1995); Ben-Hur, E. and B. Horowitz AIDS 10:1183-1190 (1996)). The most promising photochemical methods employ the use of phthalocyanines (which are activated by light in the far red (660-700 nm)) for sterilization of RBCC (Horowitz, B., et al. Transfusion 31:102-108 (1991); Ben-Hur, E., et al. J. Photochem. Photobiol. B:Biol. 13:145-152 (1992)).\nOf the phthalocyanines, silicon phthalocyanine Pc4 has been proven to be most effective for inactivation of HIV in its multiple forms (Margolis-Nunno, H., et al. Transfusion 36:743-750 (1996)). However, because Pc4 and other phthalocyanines target the lipid envelope of viruses and can, therefore, cause red cell damage, quenchers of reactive oxygen species have been used to prevent some of this damage (Rywkin, S., et al. Photochem. Photobiol. 56:463-469 (1992); Ben-Hur, E., et al. Transfusion 35:401-406 (1995)). The use of high irradiance (Ben-Hur, et al. Photochem. Photobiol. 61:190-195 (1995)) and Chremophor as the vehicle (Ben-Hur, et al. Photochem. Photobiol. 62:575-579 (1995)) also improved the specificity of viral sterilization by Pc4."}
{"text": "The invention relates generally to systems which generate images during medical and surgical procedures, and in particular, a system for generating images during medical and surgical procedures based on a scan taken prior to the procedure.\nImage guided medical and surgical procedures comprise a technology by which images, obtained either preprocedurally or intraprocedurally (i.e., prior to or during a medical or surgical procedure), are used to guide a doctor during the procedure. The recent increase in interest in this field is a direct result of the recent advances in imaging technology, especially in devices using computers to generate three dimensional images of parts of the body, such as computed tomography (CT) or magnetic resonance imaging (MRI).\nThe majority of the advances in imaging involve devices which tend to be large, encircle the body part being imaged, and are expensive. Although the images produced by these devices depict the body part under investigation with high resolution and good spatial fidelity, their cost usually precludes the dedication of a unit to the performance of procedures. Therefore, image guided surgery is usually performed using images taken preoperatively.\nThe reliance upon preoperative images has focused image guidance largely to the cranium. The skull, by encasing the brain, serves as a vessel which inhibits changes in anatomy between imaging and surgery. The skull also provides a relatively easy point of reference to which a localization system may be attached so that registration of preprocedural images to the procedural work space can be done simply at the beginning of the procedure. Registration is defined as the process of relating preprocedural images of anatomy to the surgical or medical position of the corresponding anatomy. For example, see Ser. No. 07/909,097, now U.S. Pat. No. 5,383,454, and PCT application PCT/US95/12894 filed on Oct. 5, 1995, entitled \"Surgical Navigation Systems Including Reference and Localization Frames,\" the entire disclosures of which are incorporated herein by reference.\nThis situation of rigid fixation and absence of anatomical movement between imaging and surgery is unique to the skull and intracranial contents and permits a one-to-one registration process as shown in FIG. 1. The position during a medical procedure or surgery is in registration with the preprocedural image data set because of the absence of anatomical movement from the time of the scan until the time of the procedure. In almost every other part of the body there is ample opportunity for movement which degrades the fidelity of the pre-procedural images in depicting the intraprocedural anatomy. Therefore, additional innovations are needed to bring image guidance to the rest of the body beyond the cranium.\nThe accuracy of image guided surgery is based on the identification of structures within the body that do not change shape, do not compress, nor deform between the process of imaging and surgery. Such structures are termed \"rigid bodies,\" and the bones of the skeleton satisfy this definition for a rigid body. Bones are commonly a target for medical or surgical procedures either for repair, fusion, or biopsy. Therefore, a technique is needed whereby registration can be performed between the bones or bone fragments (skeletal elements) as depicted pre-procedurally on scans and the position of these same skeletal elements as detected intraprocedurally. This technique must take into account that movement can occur between portions of the skeleton which are not rigidly joined, such as bones connected by a joint, or fragments of a broken bone."}
{"text": "The present invention relates to photonics chips and, more particularly, to structures providing optical beam steering and methods of fabricating such structures.\nA photonics chip integrates multiple optical components, such as waveguides and optical switches, and multiple electronic components, such as field-effect transistors, into a unified platform. Among other factors, layout area, cost, and operational overhead may be reduced by the integration of both types of components on the photonics chip.\nOptical beam steering finds many practical applications as, for example, part of a photonics chip. Beam steering may be accomplished by scanning with mechanical assistance through the use of movable optical elements, such as mirrors, prisms, or lenses. A distinct disadvantage of mechanically-assisted beam steering techniques is the necessity of moving the optical element with a mechanism, such as a micromechanical system, to provide directional changes needed for steering.\nImproved structures providing optical beam steering and methods of fabricating such structures are needed."}
{"text": "By wide application of wireless human input devices (such as a wireless mouse or a wireless keyboard), and even by wide application of other wireless electronic devices, in a place such as in a family or an office, there is extremely probable a problem of mutual interruption among RF frequencies; when a user finds that a wireless human input device used normally has the situation of having an incorrect action or intermittently in working and non-working, he always can not distinguish whether it is because of mutual interruption among RF frequencies or of the faults of the wireless human input device itself; such a situation becomes more and more serious, thereby providing a test system for solving the above-mentioned problem is necessary, however, the design to increase such a test system in such a conventional wireless human input device normally will derive a problem of increasing the cost of production and a problem of how to make convenience of operation of the test system etc. The problem derived normally will make a manufacturer of such a wireless human input device unwilling to study and develop the same, and users shall still suffer the above deficient situation.\nA conventional wireless human input device must be subjected to severe test before packaging and shipping in a factory, and in the last stage of packaging, if it can accept one more test, quality control of the products will get a great favor. However, conventional modes of testing mostly can only be performed with specific test tools, this indistinctly and largely increases the cost of production.\nThe inventor of the present invention provides a test method performed completely with program codes in view of the above problems to be solved pressingly of the prior arts, the test method can increase the capability of self-testing transmission quality of a wireless human input device under the condition that the hardware arrangement of the wireless human input device is maintained to be unchanged; and it is extremely convenient for performing by a user or a testing operator."}
{"text": "Golf clubs, particularly golf irons, comprise a head with a hosel and a shaft which is attached to the head by fitting the shaft into a bore in the hosel. The hosel is attached to and integral with the head. The head includes a heel, a bottom sole, a toe, a planar striking face, and a backside. The head of a conventional golf iron has a longitudinal length heel-to-toe substantially greater than the width between the striking face and the backside. On the other hand, woods generally have a relatively low length-to-width ratio. Also, the toe side of the striking face in a conventional golf iron is higher than the heel side of the striking face. This keeps the center of gravity in the center of the blade.\nGolf irons typically include a set of eleven irons, numbered one (long) through nine (short), a pitching wedge, and a sand wedge. The eleven irons of a set conventionally have varying degrees of loft angle and lie angle. The loft angle of an iron is the angle between a vertical plane which includes the shaft, and the plane of the striking face of the iron. The lie angle of an iron is the angle between the shaft and the ground (horizontal plane) when the tangent to the sole directly under the center of mass is in the horizontal plane and when the shaft lies in a vertical plane.\nThe loft angle, as the name suggests, determines how much loft is imparted to the ball when it is struck by the tilted striking face. The lie angle of the iron assures that, when the golf club is swung properly, the sole of the iron will contact the ground evenly so that the striking face will not tend to twist inwardly or outwardly.\nThe toe defines a side face on the head between the striking face and the backside. The side face in an iron extends generally back and toward the heel. The face of the toe in top view thus defines a curvilinear line between the striking face and the backside. This is known as negative draft. The outermost point on the toe (compared longitudinally with respect to the center of mass) is typically at or near the edge between the toe and the striking face. The toe in known negative draft golf irons generally extends no further away longitudinally from the center of mass than a line perpendicular to the plane of the striking face and tangent to the outermost point on the edge between the toe and the striking face. For example, the edge on some irons is beveled and this results in the outermost point being on the outside edge of the bevel.\nIn another aspect of a negative draft golf iron, a straight line approximates the face of the toe and is defined by passing a horizontal plane through the head. Two points on the plane define the line. The first point is the intersection of the plane with the toe edge of the striking face. The second point is the intersection of the plane with the toe edge of the backside. These two points in the horizontal plane define the imaginary line across the side face. The line across the side face and the plane of the striking face together form an angle which is 90.degree. or less. In known negative draft golf irons, the outermost point on the toe is at or near the edge between the toe and the striking face.\nFor example, U.S. Pat. No. 1,532,545 illustrates a golf club head with negative draft in which the length-to-width ratio is substantially high. As shown in FIG. 3 of the '545 patent, the toe face tapers back from the planar portion of the striking face toward the heel of the club. Generally, the negative draft design facilitates casting of the head and removal of the head from a mold.\nU.S. Pat. No. 2,447,967 describes a negative draft golf club head in which the length-to-width ratio is relatively low. The toe face tapers back and toward the heel at an angle of 90.degree. or less.\nU.S. Pat. No. 1,582,836 describes a metallic golf wood head with a convex toe. The head is hollow and includes an abutment behind the striking face to provide strength and to provide offset and take up the shock from hitting a golf ball. A counterweight is cast into the rear portion of the head. This head, as with solid heads, places the mass of the club primarily behind the head. The amount of mass on the toe does not significantly impact the performance of the club.\nU.S. Pat. No. 3,035,839 describes a multipurpose golf club which has a length-to-width ratio between those of conventional woods and irons. As with the wood described above, the toe is convex. The head is provided with a first sole surface on an outer toe portion and a second sole surface on a heel portion. The first sole surface intersects the striking face at a first acute angle. The second sole surface intersects the striking face at a second larger acute angle. With the second sole surface resting on the ground, the striking face has a substantial amount of lift. With the first sole surface, the pitch of the striking face is reduced. Depending on which sole surface is positioned on the ground, the club can be used for long putts and chip shots or relatively short putts. The weight of the head is concentrated toward its heel end to aid the player in making relatively long shots.\nWith many known golf clubs, it is desirable to hit a golf ball on the \"sweet spot\" or in line with the center of mass of the club. Golf balls hit with such centered shots have truer flights and travel longer distances than off-center shots. Hitting the ball slightly off-center on the striking face creates problems with the control of the direction and flight of the golf ball. The golf club tends to twist and vibrate, which absorbs the striking energy of the golf swing. The twisting and vibration impart undesirable spin to the golf ball as it leaves the striking face. Spin causes the ball to hook or to slice. The ball does not travel in the direction desired by the golfer nor does the ball travel as far as it would for a solid, on-center shot. A golf club with a larger sweet spot on the striking face helps compensate for an off-center shot by a golfer. Generally, a club with a larger sweet spot has an increased mass moment of inertia about the center of mass. The effect of an increased mass moment of inertia is a reduction of the twisting and other vibrational forms of energy loss experienced from a shot hit on the striking face at a point other than the center of mass.\nThere are several known methods of increasing the mass moment of inertia for a golf club. These methods include increasing the weight of the club and re-distributing the weight of the club. However, increasing the mass moment of inertia of a golf club head is restricted by a number of factors that relate to the playability requirements for a golf club. Increasing the weight of the club is not always possible nor desirable. For example, the total mass of the club head is determined by the shaft length, the grip weight, and the desired swing weight. Increasing the mass of the head may result in a heavier than desired swing weight. A golf club which is too heavy cannot be swung with a proper stroke. This results in off-center shots, loss of control, and shorter flights of the golf ball.\nAnother method of increasing the mass moment of inertia concentrates the majority of the club weight at the heel and the toe. Such a design reduces the golf club's tendency to twist when the golf ball is hit off-center. In a golf iron, such weight distribution is accomplished by providing a cavity in the back of the head of the golf club. The mass removed from the cavity of the club is placed around the perimeter of the backside and concentrated at the heel and the toe.\nU.S. Pat. No. 4,621,813 describes a set of golf irons having such a cavity back. The material which would otherwise be in the cavity is redistributed in the heel and toe. The redistributed material improves the resistance to twisting of the club head upon off-center impact with a golf ball. Further, these golf irons have negative draft toes.\nU.S. Pat. No. 3,059,926 describes a golf iron having an increased moment of inertia. A plug of material is added on the backside near the toe. This club also is a negative draft golf iron with the toe face forming an angle of 90.degree. or less with the striking face.\nHowever, other playability requirements also affect the distribution of mass in the golf club. These include the length and shape of the sole of the club and its contact area with the ground during a golf swing or shot. For instance, a blade with a longer length has its mass positioned farther from the center of mass. But a longer length blade has an increased sole length and this results in significantly higher ground drag. High ground drag means that the club is dragging on the ground during the shot. This may result in a slower club speed during the stroke through the golf ball and in a decreased flight distance for the golf ball. Also, a longer leading edge may inhibit the taking of a proper divot.\nThus, there exists a need in the art for a golf iron having an increased moment of inertia without the loss of desirable club dimensional and playability characteristics."}
{"text": "The present invention relates to a control system for a continuously variable transmission.\nJP-A No. 63-43837 discloses a conventional control system for a continuously variable transmission. This control system is constructed and arranged so that a target reduction ratio is achieved by a feed-back control of a shift actuator. When a motor vehicle is not driven, viz., an engine brake is applied, a feed-back gain becomes larger. This ameliorates the followability and the responsibility of shifting to the maximum reduction ratio. Thus, even if quick deceleration takes place, shifting up to the maximum gear ratio can be achieved before the motor vehicle comes to a standstill.\nA problem encountered in such a conventional control system for a continuously variable transmission is that, since the feed-back gain is kept constant regardless of change in an input torque to the transmission, viz., an engine torque, time is needed to bring a current reduction ratio to the target reduction ratio when the engine torque is large, and it is difficult to keep the target reduction ratio corresponding to the current reduction ratio even during normal operation of the motor vehicle. Specifically, with a continuously variable V-belt transmission, the reduction ratio is determined by a balance between a hydraulic pressure on a movable conical member of a pulley and a tension of a V-belt. If the engine torque increases out of a predetermined balanced reduction ratio, for example, the tension of the V-belt also increases. In this event, the hydraulic pressure on the pulley increases in a delayed manner, the reduction ratio is deviated to the maximum reduction ratio. Further, with a continuously variable traction roller transmission which is disclosed, for example, in JP-U No. 63-84451, shifting takes place in a controlling position of roller support members by the hydraulic pressure. Thus, when the engine torque increases, the hydraulic pressure within a hydraulic servo apparatus for position control of the roller support members changes in a delayed manner, resulting in an occurrence of a deviation between the current reduction ratio and the target reduction ratio. Even during normal operation of the motor vehicle, the hydraulic pressure for holding the roller support members has to increase as the engine torque becomes larger. In this event, the hydraulic servo apparatus may not be held in a predetermined position due to increased leakage from seal portions, resulting in the occurrence of a deviation of the current reduction ratio to the maximum reduction ratio."}
{"text": "The present invention relates to the recombinant DNA which encodes the N.BstNBI nicking endonuclease and modification methylase, and the production of N.BstNBI nicking endonuclease from the recombinant DNA. N.BstNBI nicking endonuclease is originally isolated from Bacillus stearothermophilus. It recognizes a simple asymmetric sequence, 5' GAGTC 3', and it cleaves only one DNA strand, 4 bases away from the 3'-end of its recognition site.\nRestriction endonucleases are enzymes that recognize and cleave specific DNA sequences. Usually there is a corresponding DNA methyltransferase that methylates and therefore protects the endogenous host DNA from the digestion of a certain restriction endonuclease. Restriction endonucleases can be classified into three groups: type I, II, and III. More than 3000 restriction endonucleases with over two hundred different specificities have been isolated from bacteria (Roberts and Macelis, Nucleic Acids Res. 26:338-350 (1998)). Type II and type IIs restriction enzymes cleave DNA at a specific position, and therefore are useful in genetic engineering and molecular cloning.\nMost restriction endonucleases catalyze double-stranded cleavage on DNA substrate via hydrolysis of two phosphodiester bonds on two DNA strands (Heitman, Genetic Engineering 15:57-107 (1993)). For example, type II enzymes, such as EcoRI and EcoRV, recognize palindromic sequences and cleave both strands symmetrically within the recognition sequence. Type IIs endonucleases recognize asymmetric DNA sequences and cleave both DNA strands outside of the recognition sequence.\nThere are some proteins in the literature which break only one DNA strand and therefore introduce a nick into the DNA molecule. Most of those proteins are involved in DNA replication, DNA repair, and other DNA-related metabolisms (Kornberg and Baker, DNA replication. 2nd edit. W. H. Freeman and Company, New York, (1992)). For example, gpII protein of bacteriophage fI recognizes and binds a very complicated sequence at the replication origin. It introduces a nick in the plus strand, which initiates rolling circle replication, and it is also involved in circularizing the plus strand to generate single-stranded circular phage DNA. (Geider et al., J. Biol. Chem. 257:6488-6493 (1982); Higashitani et al., J. Mol. Biol. 237:388-400 (1994)). Another example is the MutH protein, which is involved in DNA mismatch repair in E. coli. MutH binds at dam methylation site (GATC), where it forms a protein complex with nearby MutS which binds to a mismatch. The MutL protein facilitates this interaction and this triggers single-stranded cleavage by MutH at the 5' end of the unmethylated GATC site. The nick is then translated by an exonuclease to remove the mismatched nucleotide (Modrich, J. Biol. Chem. 264:6597-6600 (1989)).\nThe nicking enzymes mentioned above are not very useful in the laboratory for manipulating DNA due to the fact that they usually recognize long, complicated sequences and usually associate with other proteins to form protein complexes which are difficult to manufacture. Thus none of these nicking proteins are commercially available. Recently, we have found a nicking protein, N.BstNBI, from the thermophilic bacterium Bacillus stearothermophilus, which is an isoschizomer of N.BstSEI (Abdurashitov et al., Mol. Biol. (Mosk) 30:1261-1267 (1996)). Unlike gpII and MutH, N.BstNBI behaves like a restriction endonuclease. It recognizes a simple asymmetric sequence, 5' GAGTC 3', and it cleaves only one DNA strand, 4 bases away from the 3'-end of its recognition site (FIG. 1A).\nBecause N.BstNBI acts more like a restriction endonuclease, it should be useful in DNA engineering. For example, it can be used to generate a DNA substrate containing a nick at a specific position. N.BstNBI can also be used to generate DNA with gaps, long overhangs, or other structures. DNA templates containing a nick or gap are useful substrate for researchers in studying DNA replication, DNA repair and other DNA related subjects (Kornberg and Baker, DNA replication. 2nd edit. W. H. Freeman and Company, New York, (1992)). A potential application of the nicking endonuclease is its use in strand displacement amplification (SDA), which is an isothermal DNA amplification technology. SDA provides an alternative to polymerase chain reaction (PCR), and it can reach 10.sup.6 -fold amplification in 30 minutes without thermo-cycling (Walker et al., Proc. Natl. Acad. Sci. USA 89:392-396 (1992)). SDA uses a restriction enzyme to nick the DNA and a DNA polymerase to extend the 3'-OH end of the nick and displace the downstream DNA strand (Walker et al., (1992)). The SDA assay provides a simple (no temperature cycling, only incubation at 60.degree. C.) and very rapid (as short as 15 minutes) detection method and can be used to detect viral or bacterial DNA. SDA is being introduced as a diagnostic method to detect infectious agents, such as Mycobacterium tuberculosis and Chlamydia trachomatis (Walker and Linn, Clin. Chem. 42:1604-1608 (1996); Spears et al., Anal. Biochem. 247:130-137 (1997)).\nFor SDA to work, a nick has to be introduced into the DNA template by a restriction enzyme. Most restriction endonucleases make double-stranded cleavages. Therefore, modified .alpha.-thio deoxynucleotides (dNTP.alpha.S) have to be incorporated into the DNA, so that the endonuclease only cleaves the unmodified strand which is within the primer region (Walker et al., 1992). The a-thio deoxynucleotides are eight times more expensive than regular dNTPs (Pharmacia), and is not incorporated well by the Bst DNA polymerase as compared to regular deoxynucleotides (J. Aliotta, L. Higgins, and H. Kong, unpublished observation).\nAlternatively, if a nicking endonuclease is used in SDA, it will introduce a nick into the DNA template naturally. Thus the dNTP.alpha.S is no longer needed for the SDA reaction when a nicking endonuclease is being used. This idea has been tested, and the result agreed with our speculation. The target DNA can be amplified in the presence of the nicking endonuclease N.BstNBI, dNTPs, and Bst DNA polymerase.\nWith the advent of genetic engineering technology, it is now possible to clone genes and to produce the proteins that they encode in greater quantities than are obtainable by conventional purification techniques. Type II restriction-modification systems are being cloned with increasing frequency. The first cloned systems used bacteriophage infection as a means of identifying or selecting restriction endonuclease clones (EcoRII: Kosykh et al., Molec. Gen. Genet 178:717-719 (1980); HhaII: Mann et al., Gene 3:97-112 (1978); PstI: Walder et al., Proc. Nat. Acad. Sci. 78:1503-1507 (1981)). Since the presence of restriction-modification systems in bacteria enable them to resist infection by bacteriophages, cells that carry cloned restriction-modification genes can, in principle, be selectively isolated as survivors from libraries that have been exposed to phage. This method has been found, however, to have only limited value. Specifically, it has been found that cloned restriction-modification genes do not always manifest sufficient phage resistance to confer selective survival.\nAnother cloning approach involves transferring systems initially characterized as plasmid-borne into E. coli cloning plasmids (EcoRV: Bougueleret et al., Nucl. Acids Res. 12:3659-3676 (1984); PaeR7: Gingeras and Brooks, Proc. Natl. Acad. Sci. USA 80:402-406 (1983); Theriault and Roy, Gene 19:355-359 (1982); PvuII: Blumenthal et al., J. Bacteriol. 164:501-509 (1985)).\nA further approach which is being used to clone a growing number of systems involves selection for an active methylase gene (refer to U.S. Pat. No. 5,200,333 and BsuRI: Kiss et al., Nucl. Acids Res. 13:6403-6421 (1985)). Since restriction and modification genes are often closely linked, both genes can often be cloned simultaneously. This selection does not always yield a complete restriction system however, but instead yields only the methylase gene (BspRI: Szomolanyi et al., Gene 10:219-225 (1980); BcnI: Janulaitis et al, Gene 20:197-204 (1982); BsuRI: Kiss and Baldauf, Gene 21:111-119 (1983); and MspI: Walder et al., J. Biol. Chem. 258:1235-1241 (1983)).\nAnother method for cloning methylase and endonuclease genes is based on a colorimetric assay for DNA damage (see U.S. Pat. No. 5,492,823). When screening for a methylase, the plasmid library is transformed into the host E.coli strain such as AP1-200. The expression of a methylase will induce the SOS response in an E.coli strain which is McrA+, McrBC+, or Mrr+. The AP1-200 strain is temperature sensitive for the Mcr and Mrr systems and includes a lac-Z gene fused to the damage inducible dinD locus of E.coli. The detection of recombinant plasmids encoding a methylase or endonuclease gene is based on induction at the restrictive temperature of the lacZ gene. Transformants encoding methylase genes are detected on LB agar plates containing X-gal as blue colonies. (Piekarowicz et al., Nucleic Acids Res. 19:1831-1835 (1991) and Piekarowicz et al., J. Bacteriology 173:150-155 (1991)). Likewise, the E.coli strain ER1992 contains a dinD1-Lac Z fusion but is lacking the methylation dependent restriction systems McrA, McrBC and Mrr. In this system (called the \"endo-blue\" method), the endonuclease gene can be detected in the absence of it's cognate methylase when the endonuclease damages the host cell DNA, inducing the SOS response. The SOS-induced cells form deep blue colonies on LB agar plates supplemented with X-gal. (Fomenkov et al., Nucleic Acids Res. 22:2399-2403 (1994))\nSometimes the straight-forward methylase selection method fails to yield a methylase (and/or endonuclease) clone due to various obstacles, See, e.g., Lunnen et al., Gene 74(1):25-32 (1988). One potential obstacle to cloning restriction-modification genes lies in trying to introduce the endonuclease gene into a host not already protected by modification. If the methylase gene and endonuclease gene are introduced together as a single clone, the methylase must protectively modify the host DNA before the endonuclease has the opportunity to cleave it. On occasion, therefore, it might only be possible to clone the genes sequentially, methylase first then endonuclease (see U.S. Pat. No. 5,320,957).\nAnother obstacle to cloning restriction-modification systems lies in the discovery that some strains of E.coli react adversely to cytosine or adenine modification; they possess systems that destroy DNA containing methylated cytosine (Raleigh and Wilson, Proc. Natl. Acad. Sci. USA 83:9070-9074 (1986)) or methylated adenine (Heitman and Model, J. Bacteriology 196:3243-3250 (1987); Raleigh et al., Genetics 122:279-296 (1989); Waite-Rees et al., J. Bacteriology 173:5207-5219 (1991)). Cytosine-specific or adenine-specific methylase genes cannot be cloned easily into these strains, either on their own, or together with their corresponding endonuclease genes. To avoid this problem it is necessary to use mutant strains of E. coli (McrA- and McrB- or Mrr-) in which these systems are defective.\nAn additional potential difficulty is that some restriction endonuclease and methylase genes may not express in E. coli due to differences in the transcription machinery of the source organism and E. coli, such as differences in promoter and ribosome binding sites. The methylase selection technique requires that the methylase express well enough in E. coli to fully protect at least some of the plasmids carrying the gene.\nBecause purified restriction endonucleases, and to a lesser extent modification methylases, are useful tools for characterizing genes in the laboratory, there is a commercial incentive to obtain bacterial strains through recombinant DNA techniques that synthesize these enzymes in abundance. Such strains would be useful because they would simplify the task of purification as well as provide the means for production in commercially useful amounts."}
{"text": "One time programmable memory cells and corresponding memory arrays provide non-volatile memory which may be programmed once with desired values to be stored, which values then remain stored also without power being supplied. Sometimes, such one time programmable memories are also referred to as programmable read only memories (“PROM”). Such memories use configurations where elements like fuses or anti-fuses are arranged in an array and accessible by wordlines and bitlines. Fuses or anti-fuses may be irreversibly modified by applying a programming voltage, for example through breakdown of an oxide, thus programming the memory. Reading the memory then generally occurs using lower voltages than the programming voltage.\nSuch memories are for example used in telecommunication applications, sensor applications, in read-out circuits, e.g. application specific integrated circuits (ASICs) for microelectromechanical systems (MEMS) or in radio frequency switches, but are not limited to these applications. Non-volatile memories may for example be used for storing calibration data from a post-manufacturing calibration or to store software code needed for a particular device.\nIncreasingly, there is a demand for higher storage densities of such memories, i.e. a higher number of memory cells per area. Furthermore, it is desirable that production of such memories be compatible with standard manufacturing techniques, for example CMOS process compatible. Finally, high reliability of such memories is desired."}
{"text": "1. Field of the Invention\nThe invention relates in general to a vertical channel transistor structure and a manufacturing method thereof, and more particularly to a vertical channel transistor structure with a narrow channel and a manufacturing method therefor.\n2. Description of the Related Art\nAlong with other advances in semiconductor manufacturing technology, the resolution of current semiconductor elements has reached the nanometer level. For example, the reduction in gate length and element pitch in memory units is carried on continually. Although the technology of photolithography has improved greatly, currently manufactured planar transistor structures have reached the limit of resolution, and the transistor elements manufactured thereby are apt to have the problems of electrostatic discharge (ESD), leakage, and decrease in electron mobility, resulting in short channel effect and drain induced barrier lowering (DIBL) effect. Thus, the double-gate or tri-gate vertical channel transistors capable of providing higher packing density, better carrier transport and device scalability, such as the fin field effect transistor (fin FET) for instance, have become transistor structures with great potential.\nThe fin FET transistor has a vertical channel that can be formed on the two lateral surfaces of the fin FET transistor and turns on the current by the double-gate or the tri-gate, hence having higher efficiency than conventional planar channel transistors.\nWhen manufacturing a fin FET element with high resolution, expensive processes such as the photolithography process and the E-beam process are required. Therefore, the throughput can hardly be increased and large-scale production is difficult to achieve. There is another manufacturing method which reduces the channel width by applying oxidation to the etched channel. However, the element formed according to the above method has poor uniformity and unstable quality."}
{"text": "L-cysteine is a sulfur-containing amino acid that can be synthesized de novo from methionine and serine in adult humans. L-cysteine performs a variety of metabolic functions. For example, L-cysteine is involved in growth and protein synthesis and it is a precursor for glutathione, an important intracellular antioxidant.\nL-cysteine is generally classified as a non-essential amino acid or “semi-essential” amino acid because it can be synthesized in small amounts by the human body. However, some adults can still benefit from L-cysteine supplementation. Further, L-cysteine has been classified as conditionally essential in some cases. For example, L-cysteine can be conditionally essential in preterm infants due to biochemical immaturity of the enzyme cystathionase that is involved in L-cysteine synthesis. Thus, there are a number of circumstances in which L-cysteine supplementation can be desirable.\nThe subject matter described herein addresses the shortcomings of the art by providing L-cysteine compositions that facilitate the desired supplementation but with an exceptional safety, purity and stability profile."}
{"text": "When the strings of an electric guitar are plucked, sound waves are generated in the air and electrical signals are also created in an associated electrical circuit. The vibrations of the strings, which are made of a magnetically permeable metallic material, induce signals in a device known as a pickup. The pickup is positioned underneath the strings of the instrument in the upper plate. The pickup operates on an electromagnetic principle and provides a signal to an output circuit.\nThe most essential components of a guitar pickup are a permanent magnet and a coil of wire. The magnet generates a magnetic field that passes through the pickup coil and also extends into a space occupied by at least one string of the instrument. Vibration of the string changes the reluctance of the magnetic path and creates disturbances in the magnetic field proportional to the string vibration. The changing magnetic field in the pickup coil in turn induces an electrical signal in the coil. From the output of the pickup a circuit connection is made to an amplifier and speaker system. Impedance matching in the electrical system may be required to transfer the music signal effectively.\nThe energy level of the string vibration is very small and design of the pickup is calculated to obtain maximum signal output. Pickups often include a pole piece in addition to the permanent magnet so as to concentrate the magnetic field where it is most needed and thus maximize the output signal.\nIn a musical performance the direct sound output of the electric guitar--i. e., its acoustical output--is so low that it does not add appreciably to sound output originating from a pickup or pickups that is disseminated through an amplifier and loudspeakers. In some instances an artist may wish to make a recording using only the sound output developed in the electrical system. Regardless of which way a performance or recording is to be handled, the fact remains that a high quality of musical sound output from each pickup is required.\nThe magnetic field of the pickup necessarily constitutes a load on an associated string, causing some damping or diminution in the string vibration. As a result, there may be some non-linearity with consequent distortion and loss of quality in both the acoustical sound output and the electrical sound output. Guitar musicians may differ, however, in their evaluation of musical quality, and in some instances a certain amount of distortion of a pure musical sound may be preferred.\nAs is well known, string vibrations occur at multiple frequencies--that is, a string that is vibrating at a certain fundamental frequency is concurrently vibrating at two, three, four, etc., times that fundamental frequency. At these higher overtone frequencies the length of each vibrating string segment becomes quite small. Thus, the vibrations of a string at different locations along its length will be out of phase with each other for some frequencies. If a dual pickup coil detects different portions of the string vibration that are out-of-phase with each other, the result will be a significant loss of signal strength at the corresponding frequencies.\nThe guitar is an instrument with six strings, each having its own assigned frequency. Theory might indicate that there should be a separate pickup for each string. However, practical cost considerations have resulted in the practice of using an elongated coil that can pick up vibration signals from all of the strings at once. The coil is made long enough to span all six of the strings and is positioned with its longitudinal axis essentially perpendicular to the length of the strings.\nIn every pickup coil a certain number of turns are needed in order to effectively respond to the changes in the magnetic field; i.e., to produce a useful output signal in response to the string vibration. Providing a coil with the desired number of turns necessarily results in inductance, resistance, and capacitance. For a given wire size the resistance of the coil is essentially proportional to the length of the wire and hence to the number of turns in the coil. The same is true for the capacitance. Changes in wire size and number of turns may result in a change in the musical character of the sound output that would not be acceptable to musicians using the instruments.\nA well-known problem of the electric guitar is that the pickup coil, in addition to its desired function of picking up string vibrations, also tends to pick up electrical noise and interference signals from various extraneous sources, such as power circuits, radio and television equipment, and the like. It has long been well known to utilize a two-coil pickup in which the coils are interconnected in such a way as to balance out the extraneous signals. Such pickups are known in the trade as \"humbuckers\".\nThe operating principle of the humbucker is that the two coils are connected in opposite electrical polarities so that the noise signals which are electrically induced in them will cancel each other out. At the same time, the magnetic circuits of the humbucker's two coils are so arranged that the signals magnetically induced in them from the string vibrations will be added together rather than being cancelled. An important requirement for a humbucker is that the electrical impedances of the two coils must be substantially identical. Otherwise the noise signals induced in the two coils will be unequal and will not completely cancel out."}
{"text": "present invention relates to a micropin array and the method of producing the same, and more specifically relates to a set of I/O pins for use as an external connector of a PGA package etc. and the method of producing the same.\nConventionally, in the PGA package of an IC device, there are utilized I/O pins in the form of a set of micropins which are separately and independently fixed to an LSI chip. Such micropins are produced such that a lead wire material is treated by mechanical working such as stamping with header processing to form individual micropins. Namely, the micropins are produced in manner similar to the production method of the typical I/O pins attached to the typical PGA package. There may be other methods such as utilizing photolithographic technology to form chemically the micropins. These methods of producing micropins are disclosed, for example, in Japanese Patent Publication Nos. 62-24916, 62-32591, 62-32592 and 63-28515.\nIn use of the conventional micropins for electrical connection to an LSI chip, the micropins are individually fixed to the LSI chap in a given alignment pitch, thereby disadvantageously requiring a highly accurate alignment tool. Further, since the micropins are connected to the LSI chip one by one, micropins are easily curved or folded to thereby disadvantageously degrade the mechanical strength of an assembly.\nIn the conventional method of producing micropins according to the mechanical working, there is practical limitation in diameter and length dimension of micropins due to dimensional error in the header processing. Therefore, it is practically difficult to produce a pin having a diameter less than 0.1 mm.\nOn the other hand, a fine pin having a diameter less than 0.1 mm can be formed by the photolithographic technology. However, the aspect ratio thereof, i.e., pin length/pin diameter ratio cannot be increased by this method."}
{"text": "This invention relates to a self-locking nut of the anti-reversing type, to a method for fabricating such a self-locking nut, and to such method including the new use of a material suitable for the anti-reversing pin component of the self-locking nut.\nThis invention is further concerned with self-locking nuts of a type which are suitable for use in the fabrication of structures such as microwave towers and highway structures. In the fabrication of such structures it is desirable to use nuts and bolts which have a very effective protective coating to minimize deterioration due to rust, corrosion, etc.\nA very effective protective coating for steel nuts and bolts is referred to as \"hot dip galvanized\"; and one characteristic of this protective coating is that the thickness of the coating is not uniform. With that characteristic, many known types of lock nuts may not be used effectively, since where the hot dip coating is particularly thick, the binding may be so severe that the self-locking nut may not be turned.\nA form of self-lacking nut which has been found effective for use with bolts treated by the hot dip galvanized process, is a nut having an anti-reversing pin which projects generally radially inward from the nut for interference engagement with the threads of the bolt. In the manufacture of such anti-reversing nut, the interference pin projects radially into the thread bore; and in use that pin will be deflected in one direction or the other, about 30.degree. from the radial position for example, to assume an anti-reversing position relative to the direction in which the nut is being turned on the bolt. The tip of this pin is configured to be received partially within the groove of the bolt thread and to bite into the bolt material when urged in a reverse direction. The pin, which is anchored in the nut, is necessarily elastic to maintain the interference engagement with the bolt thread and maintain the anti-reversing force, which enables it to perform its function as a self-locking nut.\nWhen it is desired to remove the nut from the bolt, a high torque must be applied which effects the swinging of the anti-reversing pin and accompanying deflection in the opposite direction from the radial position, whereby the nut then becomes anti-reversing in the opposite direction.\nTo provide an effective self-locking nut for a hot dip galvanized bolt, it is desirable that the nut be relatively free spinning, either while threading the nut onto the bolt or while unthreading the nut from the bolt, after overcoming the high anti-reversing torque, for several reasons. One reason is to minimize any galling or other defacing of the protective coating on the threads which may result in exposure of the base metal and thereby subject the bolt to rust or other corrosion. Another reason is to facilitate the construction of structures being asembled with nuts and bolts and to promote the safety of the workers who may be laboring at either high or otherwise hazardous locations. When working at high locations, such as on microwave towers, the use of power tools is sometimes impractical; and a requirement for applying high torque to these fasteners may be hazardous to the workers who must at all times be concerned with body balance and stability.\nWhile self-locking nuts having inwardly directed interference pins as above described are known in the art, a problem with some of these nuts is that the pins break off after being deflected because of lack of sufficient elasticity. Those nuts of course must be discarded; and the reliability of the nuts comes into question. That is, do the anti-reversing pins have sufficient elasticity to perform the intended anti-reversing or locking function.\nAn object of this invention is to provide an improved self-locking nut of the anti-reversing type, and a method for manufacturing same.\nAnother object of this invention is to provide an improved self-locking nut of the anti-reversing type which is relatively free spinning while running the nut onto or off from a bolt, and a method for manufacturing same.\nA further object of this invention is to provide a self-locking nut of the anti-reversing type which is adapted to be used repeatedly, and a method for manufacturing same.\nStill another object of this invention is to provide an improved self-locking nut of the anti-reversing type having an anti-reversing pin which projects into interference engagement with the bolt and which is adapted to be oscillated between alternative anti-reversing positions, and a method for manufacturing same.\nA still fruther object of this invention is to provide a self-locking nut of the anti-reversing type as set forth in the preceding object, wherein the anti-rotation pin is fabricated from quarter hard stainless steel welding rod.\nAnother object of this invention is to provide an improved self-locking nut of the anti-reversing type having an anti-reversing pin which projects into interference engagement with the bolt, and which has the required elasticity for repeated oscillations between alternative anti-reversing positions, and a method for manufacturing same.\nThe novel features and the advantages of the invention, as well as additional objects thereof, will be understood more fully from the following description when read in connection with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to an air-inlet assistor for a vehicle engine. More particularly, the present invention is directed to a system for automatically keeping the air supply for an engine balanced.\n2. Prior Art\nIn accordance with the operational principles of conventional vehicle engines, and with reference to FIG. 4, the amount of fresh air from the outside passing through the air cleaner 53 for supplying the running engine is dependent on the drawing in force of the inlet manifold 52 through the carburetor 51. The inlet pipe 54 is connected to the base stand 55 that is attached to the air cleaner 63, the inlet pipe being connected through an inlet valve 56 and a back valve 57, to form an air-inlet assistor for improving the full burning of fuel in the engine. The base stand 55 has a cubic shape and there is an L or S shaped channel passing through it, with one end connecting to the air cleaner 53 for cleaning the air, and the other end being connected to the inlet valve 56. The outlet of the back valve 57 is connected to the carburetor 51 of the engine 5. So when starting the engine, the drawing in force created in the inlet manifold 52 can draw the air into the engine 5 for burning the fuel.\nAccordingly, the above-mentioned air-inlet controlling device has the following shortcomings:\n1. The conventional air-inlet assistor cannot be adjusted automatically. If the user finds the air-inlet not providing sufficient air, so as to influence running, he would drive to the garage and leave it to a technician. Therefore, it is not convenient for the user. PA1 2. Adjustment of the conventional device is by manually changing the opening of the inlet valve, causing an additional factor that may influence the intake of air. It too much air is drawn, a mono-cylinder effect can be generated, resulting in a shaking of the vehicle. If the adjustment is in the opposite direction, the burning of fuel in the engine may be incomplete, increasing the waste of gas and adding to atmospheric pollution. PA1 3. Due to the many shortcomings of the conventional device, its production costs are increased. Also, special technicians are needed to install it, making it not practical for high volume production.\nAn object of the present invention is to provide an air-inlet assistor that can automatically adjust the air-inlet accurately, according to the vehicle, the engine, the differences in the drawing in force from the inlet manifold of the engine, and the different working state of the engine, to avoid the intake of air that is greater or less than needed, thereby keeping the engine running smoothly."}
{"text": "Polymeric film structures and packages made therefrom are useful in the packaging field for the packaging of food products, especially respiring food products such as natural cheeses. These film structures and the packages made therefrom generally contain multiple layers of polymers in which each layer adds certain physical or chemical properties to the completed film or package made therefrom.\nIn the packaging of respiring food products such as natural cheese (i.e., swiss-type cheese) certain packaging problems exist. These packaging problems exist due to process by which the respiring product is made and because of the packaging requirements of the final product. For example swiss-type cheeses are made utilizing specific molds or bacteria to produce the “eyes” which are characteristic of this type of cheese. Specifically, swiss-type cheeses are ripened by typically adding bacteria such as Propionibacter Shermanii to form the “eyes” of the cheese. These “eyes” are formed as gas pockets of CO2 which is given off by the swiss-type cheese. This CO2 elimination not only occurs during production of cheese but continuing during the “life” of the product including the packaged product. Too much CO2 inside the package causes the “package to “pillow.” “Pillowed” packages are negatively received by the consumer. Also, natural food products such as swiss-type cheeses are affected by atmospheric oxygen (O2) during the transporting and storing of this type of product in a package. If the permeability of O2 is too rapid, the product “life” is shorter. Shorter product life affects the financial aspects of the product. Therefore, these inherent problems which are associated with respiring food products such as swiss-type cheeses must be addressed by utilizing film structures and packages made therefrom which will deal with these problems in an efficient manner. In addition to addressing the above problems, the film structures and packages must also provide stiffness, moisture barrier and maximize flex crack resistance. Also, the film structures and packages must be cost-effective.”"}
{"text": "Skid steer loader machines are manufactured to provide a compact work machine that is maneuverable and durable in various applications. Skid steer loader machines are propelled by independently driven wheels on opposite sides of the vehicle that can be driven at different speeds and in different directions. Some skid steer loader machines use a loader body or lower frame assembly, in conjunction with a subframe or upper frame assembly, to define a main frame assembly for the loader machine. The lower frame assembly typically includes axle openings within which axles are disposed. The main frame may be assembled in a variety of ways to provide effective serviceability and durability.\nTracked skid steer loader machines or multi-terrain loader machines utilize an endless track undercarriage in place of the independently driven wheels to provide greater traction in mud, snow, ice and the like. An all-terrain vehicle is disclosed in U.S. Pat. No. 3,728,909 issued on Apr. 24, 1973 to Herbert A. Jesperson that is designed for high capability of travel over difficult and rough ground contours, conditions and textures. In particular, the Jesperson patent includes a body that has a center section with a bottom wall and front and rear walls extending from the bottom wall. Opposing side walls extending from the bottom wall and from the front and rear walls generally complete the center section. Located transversely on the underside of the bottom wall are a plurality of spaced parallel downwardly open channels which are formed in the continuous bottom wall surface. The channels cooperate with the track suspension group to permit modular assembly while, at the same time, affording maximum vertical clearance under the bottom wall and between the tracks, consistent with maintaining water-tight integrity of the body and maximum depth to afford sufficient space in the center section for the power group and the user's legs. In this regard, the track suspension group includes a frame having a plurality of structural cross members which are respectively and individually located in the respective downwardly open channels. The cross members are connected to the body by means in the form of plates which are welded to the cross members and adapted to span the ends of the channels and contact the exterior surface of the side margins around the channel ends. Screws or bolts are used to attach the plates to the side walls of the body. Although the open channels provide vertical clearance and depth, axial loads from the upper portion of the machine are directed on the screws or bolts attaching the cross members with the body. Therefore, such screws or bolts are likely to incur shear or torsional stresses that may damage or destroy the screws or bolts and, thereby, weaken the assembly. Further, the depth utilized to mount the cross members is not sufficient to lower the center of gravity of the machine to a maximum level.\nThe present invention is directed to overcoming the problems, as set forth above, by providing a structural design that enables a multi-terrain loader machine to be modularly assembled by mounting a solitary lower frame assembly or a fully assembled main frame to a pre-assembled undercarriage in a manner that distributes axial loading on the undercarriage."}
{"text": "Detergents are well-known in the art. The detergent compositions that are currently available are suitable for their intended purposes, i.e., effectively cleaning, and leaving previously soiled fabrics, clothing, dishes, or utensils in a generally spot-free, clean condition.\nDetergents can include dishwashing detergents, carpet cleaning detergents, and others. These and other detergents contain some combination of one or more of three ingredients, including bleach, caustic soda, and phosphates. These substances can be deleterious, for various reasons.\nFor example, phosphates are minerals that act as water softeners, and are considered by some to be among the worst pollutants found in detergents. Phosphates are a nutrient, and act as a fertilizer for algae. Thus, when phosphates enter waterways, they promote the growth of algae and other plants. In the presence of large amounts of phosphates and other similar nutrients, excessive algae growth occurs. This causes odors and creates hypoxic conditions.\nSome states have banned the use of phosphates in all detergents, other than automatic dishwasher detergents. Thus, phosphate-free general purpose detergents and clothing detergents are readily available.\nCaustic soda is also considered to be a contaminant of groundwater. For example, alumina plants use and discharge caustic soda into the groundwater adjacent their plants, and such groundwater contamination is considered to be a significant environmental problem. Moreover, the caustic soda can impart to an automatic dishwashing detergent a noticeable odor, which can be described as a “choking-type” odor. Moreover, an etched film can be irreversibly formed on the surface of glassware that has been repeatedly washed with automatic dishwashing detergents which contain caustic soda.\nBleach provides dishwashing detergents with some disinfecting capabilities. While bleach in groundwater is not generally deemed to be deleterious, the removal of bleach from an automatic dishwashing detergent would lower the cost of that dishwashing detergent. Bleach, especially when combined with caustic soda, can contribute to the choking-type odor that may be perceived in some other automatic dishwashing detergent products.\nDishwashing detergents are not the only kinds of detergents that exhibit certain distinctive deficiencies. Conventional carpet cleaning detergents can also exhibit certain deficiencies. In particular, conventional carpet cleaning detergents can have certain deleterious effects, and can have relatively high foaming characteristics.\nFirst among these deficiencies, some conventional carpet cleaning detergents can have an alkaline pH. Such alkaline pH detergents can damage the mechanical parts of automatic carpet cleaning machines. In contrast, a carpet cleaning detergent with a neutral pH (7.0) would not damage the mechanical parts of such machines.\nHigh foaming in a carpet cleaning detergent is relatively disadvantageous, as well. A “foamy” detergent will sit atop the surface of the carpet, and will thus have a tendency to not enter the sub-surface fibers of the carpet. As a result, a larger amount of such a high-foaming detergent would be necessary, in order to obtain the same results as one could obtain by using a smaller amount of a low-foaming detergent.\nIt would be desirable to provide detergents that are free of phosphates. It would further be desirable to provide a detergent which is free of caustic soda. Finally, it would be desirable to provide a detergent which is free of bleach."}
{"text": "The present invention relates to a method of and a system for recording and reproducing information data, which may be suitably used for protecting one\"\"s copyright when music and image information data or the like are transmitted for distribution by way of electronic communication system.\nRecently, audio music data, video image data and computer program data are usually transmitted for distribution by way of electronic communication system. Such an information distribution is carried out under a prerequisite that the information data can only be distributed by a legally allowed person or company. Namely, in order to forbid an unlawful copy so as to avoid an economic damage to an original author, there have been a regulation prescribing that once the above information data has been received by one recording-reproducing apparatus, the same information data must not be copied from one recording-reproducing apparatus to another. Here, xe2x80x9ccopyxe2x80x9d means that the same identical data is existing in two or more mediums.\nHowever, under the above-mentioned regulation, there has been the following problem which has been experienced by users who have legally obtained the above-described information data.\nFor example, if a person wants to enjoy the music information recorded in an internal memory means such as a hard disc of a computer, he or she has to go to the place where the computer has been set. Alternatively, the computer has to be transported to a place where the person is.\nMoreover, under the above regulation, the above-described information data must not be copied to an easily portable machine such as an MD player.\nIt is an object of the present invention to provide an improved method and an improved system for recording and reproducing desired information data, such that when information data has already been transmitted from one recording medium to another, information reproduction from the one recording medium may be forbidden (if necessary), thereby preventing the same identical information from being recorded on several different recording mediums, thus effectively preventing an unauthorized copying.\nIt is another object of the present invention to provide an improved method and an improved system for recording and reproducing information, such that the unit of information data to be transmitted is considerably small so that it is possible to almost simultaneously reproduce information data on several different apparatus, without the necessity of duplicating these informations.\nAccording to the present invention, there is provided an information reading and recording method for effecting data transmission among a plurality of information processing apparatus, said method comprising the steps of: reading and transmitting information data recorded on a recording medium of one information processing apparatus, in accordance with a transmission request from another information processing apparatus; rendering ineffective a predetermined unit amount of information data recorded on the recording medium of said one information processing apparatus, at a time a predetermined unit amount of information data has already been transmitted therefrom; recording the predetermined unit amount of transmitted information data in said another information processing apparatus which has received the transmitted information data.\nIn one aspect of the present invention, said another information processing apparatus is adapted to record the transmitted information data while at the same time reproducing the transmitted information.\nIn another aspect of the present invention, each predetermined unit amount of information data is provided with a transmitted code indicating that said unit amount of information has already been transmitted therefrom, so that it is allowed to determine whether the information to be transmitted should be reproduced.\nFurther, according to the present invention, there is provided an information reading system for reading information data from a recording medium and transmitting the information data in accordance with a transmission request, said system comprising; reading means for reading-out information data from the recording medium in accordance with a control signal; transmitting means for transmitting the read-out information data; storing means for storing transmitted code indicating that information data has already been transmitted, corresponding to each unit amount of information to be transmitted; detector means for detecting transmitted code; and control means for producing control signal for controlling the reading of said information data, and for producing another control signal for forbidding the reading of the information data when a corresponding transmitted code has been detected.\nIn another aspect of the present invention, the transmitting means for transmitting the read-out information comprises; confirmation means for confirming whether or not a plurality of different information processing apparatus are allowed to be connected with one another for exchanging predetermined data in accordance with a predetermined protocol; and an allowance issuing means for issuing an allowance for transmitting information data when it has been confirmed by the confirmation means that the plurality of different information processing apparatus are all legally connected with one another for exchanging predetermined data in accordance with a predetermined protocol.\nMoreover, according to the present invention, there is provided an information recording system adapted to record information containing transmitted code corresponding to information data to be transmitted, said system comprising: determining means for determining whether or not there is existing transmitted code corresponding to transmitted information data; and recording means for recording the transmitted information data when it is determined by the determining means that there is not existing said transmitted code, and for erasing the transmitted code when it is determined that there is existing said transmitted code.\nIn a still further aspect of the present invention, said recording means includes: confirmation means for confirming whether or not a plurality of different information processing apparatus are allowed to be connected with one another for exchanging predetermined data in accordance with a predetermined protocol; and an allowance issuing means for issuing an allowance for transmitting information data when it has been confirmed by the confirmation means that the plurality of different information processing apparatus are all legally connected with one another for exchanging predetermined data in accordance with a predetermined protocol. With the use of the present invention, when information data is transmitted from a recording medium of one apparatus to a recording medium of another apparatus, it is allowed to forbid the reading of information data from the recording medium of the one apparatus, thereby ensuring that the same identical data information will not be existing in several different recording mediums. In this way, information data can be recorded in a recording medium such as an MD (Mini Disc) having an excellent portability and owned by a person of good will, without damaging a concept of duplicating information data. Further, with regard to the provision of transmitted flag, since the information data to be transmitted has been processed to be present in quite small unit such as ECC block unit, information data such as music data can be reproduced almost simultaneously in a plurality of apparatus (among which information data is transmitted from one another), without duplication of the information data. Furthermore, since the movement of information data is administrated using data transmission unit and transmitted code, it is allowed not to erase the information data in transmission unit, thereby making it possible to reduce an overhead when the movement of information data is repeated. In addition, since the information data is transmitted, read and recorded under a condition where a mutual confirmation has been completed, it is sure to prevent any unlawful duplication of the information data and to prevent the information data from being accidentally erased."}
{"text": "This invention relates to a process for purifying a highly unsaturated fatty acid or its ester. More particularly, the invention relates to a process for purifying eicosapentaenoic acid or its ester from a mixture containing the highly unsaturated fatty acid or its ester, by a silica gel column chromatography.\nHighly unsaturated fatty acids or their esters such as eicosapentaenoic acid (hereafter called xe2x80x9cEPAxe2x80x9d) and docosahexaenoic acid (hereafter called xe2x80x9cDHAxe2x80x9d) contained in fish oil such as sardine oil have been used in the field of drugs and health foods, since they have a pharmacological effect in vivo and are effective against diseases such as arteriosclerosis obliterans and hyperlipemia.\nKnown processes for the purification of these unsaturated fatty acids include a urea-addition process, a precision distillation process, a column chromatography process and a supercritical fluid extraction process. These processes have been used singly or in combination for purifying highly unsaturated fatty acids.\nHowever, the urea-addition process is not satisfactory as a purification process for removing eicosatetraenoic acid (hereafter called xe2x80x9cETAxe2x80x9d) or its ester, because of its being poor in urea-adduct forming ability. The precision distillation process needs an operation at high temperature, and may easily cause denaturation by polymerization or isomerization. In the purification of EPA or its ester, the process has a great difficulty in removing ETA or its ester, because ETA and EPA have the same carbon numbers and make a small difference in boiling point. Further, since the supercritical fluid extraction process deals with high-pressure fluid, it has a difficulty in installation of equipment and it is unsuitable for mass separation and purification on an industrial scale.\nAs a process for the purification of highly unsaturated fatty acids using the column chromatography method, Japanese Patent Kokai 5-222392 discloses a combination of precision distillation and reverse phase partition chromatography using a silica gel on which an octadecyl group is held. Further, Japanese Patent Kokai 9-151390 discloses a purification process using a carrier having a silver salt carried on silica gel. These purification processes can separate and purify EPA and its ester from mixtures of highly unsaturated fatty acids. The former process is excellent in the separation and purification, but it has the problem that the cost of filler becomes high because the octadecyl group is held on the silica gel surface. The latter process has the problem that the cost is high due to use of the silver salt, that the method for the preparation of the carrier carrying the silver salt is complicated, that the extraction operation is complicated with repeated agitation and filtration, and that the product may be contaminated with a silver compound by using a particular solvent.\nUsually, ETA and its ester are contained in the raw material of EPA and its ester. However, ETA and its ester are impurities very difficult to separate and remove in the purification of EPA and its ester, as mentioned above.\nThus, there has been a demand for a process for the purification of eicosapentaenoic acid or its ester, by which the problems encountered in the prior art can be overcome and the contents of ETA and its ester can be remarkably reduced by selectively removing impurities.\nAccordingly, an object of the present invention is to provide a process for selectively separating and purifying eicosapentaenoic acid or its ester effectively at low cost.\nA further object of the present invention is to provide a process for selectively separating and removing ETA or its ester by a normal phase column chromatography using the silica gel under a specific condition.\nThe present invention provides a process for separating and purifying eicosapentaenoic acid or its ester from a mixture containing a highly unsaturated fatty acid or its ester, which comprises developing the mixture containing the highly unsaturated fatty acid or its ester with a mixed solvent of an ether solvent and a hydrocarbon solvent, using a column filled with a silica gel having a particle diameter of 1 to 100 micrometers.\nIn this process, the silica gel used as a filler of the column is required to have a particle diameter in the range of 1-100 micrometers. If the silica gel of small particle diameter is used, pressure loss will become large, but the amount of raw material loaded and the purifying capability will be increased due to its increased surface area. On the other hand, if the silica gel of large particle diameter is used, pressure loss will become small, but the amount of raw material loaded and the purifying capability will be reduced. Thus, it is required to choose the silica gel of suitable particle diameter in view of the amount of raw material loaded and purifying capability. For the selective separation of ETA or its ester, it is preferable to use the silica gel providing small pressure loss and having as large a specific surface area as possible. In the present invention, the form of silica gel may be spherical or fragment as far as the silica gel has the above-defined particle diameter, but spherical form is preferable. In addition, it is preferable to use porous silica gel having a pore size of not less than 4xc3x9710xe2x88x923 micrometers.\nThe amount of this silica gel used varies depending on the quantity of raw material to be processed and/or size of a column, etc., but it is usually within the range of 5-50 grams, preferably 10-40 grams per gram of raw material, i.e., the mixture comprising the highly unsaturated fatty acid or its ester.\nIn the present process, the mixture containing the highly unsaturated fatty acid or its ester is charged in the column filled with silica gel, and then developed with an organic solvent serving as a developing solvent. As the developing solvent, a mixed solvent with the highest purifying capability can be selected by varying the polarity of solvent. According to the present invention, the mixed solvent of a hydrocarbon solvent and an ether solvent is advantageous for the separation and purification of eicosapentaenoic acid or its ester.\nThe hydrocarbon solvents in the mixed solvent can include, but are not limited to, hydrocarbons of 5-8 carbons, preferably alkane hydrocarbons, for example, n-pentane, n-hexane, n-heptane and n-octane. The ether solvents can include, but are not limited to, dialkyl ethers, for example, diisopropyl ether, diethyl ether and methyl tert-butyl ether. In the present process, the developing solvent can be used in any combination of the hydrocarbon solvent and the ether solvent, but a mixed solvent of n-hexane (hereafter called xe2x80x9cn-Hexxe2x80x9d) and diisopropyl ether (hereafter called xe2x80x9cIPExe2x80x9d) is especially preferable. A mixing ratio of these solvents can be suitably decided so as to provide optimum conditions for the purification of the desired EPA or its ester. For a mixed solvent of n-Hex and IPE, for example, the mixing ratio is a volume ratio of n-Hex:IPE=90-99:10-1, preferably 95-99:5-1.\nIn the preferred embodiment of the present invention, eicosapentaenoic acid or its ester can be separated from the mixture of the highly unsaturated fatty acid or its ester by a column chromatography, which comprises developing said mixture with a mixed solvent of n-Hex and IPE as a developing solvent, using a column filled with silica gel having a particle diameter of 1-100 micrometers, at a linear velocity of 1-10 cm/min.\nFrom another viewpoint, the present invention relates to a process for selectively removing ETA or its ester from a mixture comprising a highly unsaturated fatty acid or its ester containing eicosatetraenoic acid (ETA) or its ester, by the above-described column chromatography.\nThe column used in the process of the present invention is prepared by suspending silica gel in an organic solvent to slurry it, filling the slurried silica gel in a medium pressure column and flowing an organic solvent about 1-10 times the volume of the column to stabilize the column. A preferable organic solvent used when filling silica gel in slurry is acetone. It is preferable that an organic solvent used in the stabilization of column is the same as the developing solvent used in the purification.\nThe optimum amount of the mixture containing the highly unsaturated fatty acid or its ester which should be charged in the column can be varied depending on the composition ratio of each fatty acid in the raw material, but it is in the range of 2 to 20% by weight, preferably 3 to 15% by weight, based on the weight of silica gel in the column. Preferably, the raw material, i.e., the mixture containing the highly unsaturated fatty acid or its ester is dissolved in a developing solvent at most 5 times, e.g., twice the weight of the mixture, and thereafter this solution is flowed through the column.\nFor the column used in the process of the invention, preferable is a medium pressure column with the column length diameter ratio of 2 or more and the theoretical plate number of 15000-20000 plates/m. The pressure resistance of column may be in the range of 1-6 MPa, depending on the length and the linear velocity of column. The column is preferably operated at 10-30xc2x0 C., e.g., at room temperature. In the separation of ETA or its ester, it is preferable that the mixture of raw material is developed while applying the pressure so as to provide a linear velocity of 1-10 cm/min in the column.\nThe mixture containing the highly unsaturated fatty acid or its ester used as a raw material in the present process can be obtained from fish oils extracted from mackerels, sardines, codfishes, etc., but is not limited especially. This raw material mixture contains ETA, EPA, DHA and other components such as fatty acids of less than 20 carbons.\nThe invention is further illustrated by the following non-limiting examples and comparative examples."}
{"text": "The present invention relates to electronic circuits and to digital communication. More particularly, the present invention relates to an advanced method for handling data units (packets) within a communication device by efficiently determining a flow identification for a data unit using a simplified and protocol independent method.\nThe present invention has a number of applications in advanced telecommunication systems and networks. One envisioned application of the invention is in a state-of-the-art router (or layer three switch) for handling data packets. These data packets may be constructed according to a number of well-know packet transmission protocols in a layered protocol suite. Routers are packet switching devices characterized by their ability to perform intelligent routing of packets on a per-packet basis. Routers generally can handle packets or frames of data formatted according to a number of different protocols. A layer three switch is a type of router that performs layer three routing of packets largely or exclusively using dedicated logic circuitry.\nOther applications of the invention include fast recursive look-up for grouping any kind of discreet data unit quickly into groups for group-based handling.\nExplanation of both the prior art and the invention will be better understood with reference to specific examples, however this should not be taken to limit the invention to the particular examples described. The invention is generalizable to other similar types of communication devices or to other computing devices handling data is ways that are analogous to data handling in modern networks. The invention should therefore not be limited except as provided in the attached claims.\nNetworking Devices Standards\nThis specification presumes familiarity with the general concepts, protocols, and devices currently used in LAN networking and WAN inter networking applications such as, for example, the IEEE 802 and ISO 8802 protocol suites and other series of documents released by the Internet Engineering Task Force that are publicly available. LANs are arrangements of various hardware and software elements that operate together to allow a number of digital devices to exchange data within the LAN and also may include internet connections to external wide area networks (WANs).\nPackets\nIn the most common LANs and WANs in use today, data is transmitted among network devices and processed within devices as series of independent packets, each packet containing at least one header having at least a destination address specifying an ultimate destination and generally also having a source address. Headers may also include other handling information such as an indication of the different protocols used to format a packet, a transmission priority, etc.\nFIG. 1 depicts an example of a packet as it may be transmitted on a network. The example shown is essentially an Ethernet packet, having an Ethernet header 202 and a 48-bit Ethernet address (such as 00:85:8C:13:AA) 204, and an Ethernet trailer 230. Within the Ethernet packet 200 is contained, or encapsulated, an IP packet, represented by IP header 212, containing a 32 bit IP address 214 (such as 199.22.120.33). Packet 200 contains a data payload 220 which holds the data the user is interested in receiving or holds a control message used for configuring the network. Many other types and configurations of packets are presently known in the networking art and it is assumed that additional packet protocols will be developed in the future.\nLayers\nAn additional background concept important to understanding network communications is the concept of layered network protocols. Modern communication standards, such as the TCP/IP Suite and the IEEE 802 standards, organize the tasks necessary for data communication into layers. At different layers, data is viewed and organized differently, different protocols are followed, and different physical devices may handle the data traffic. FIG. 2 illustrates one example of a layered network standard having a number of layers, the Physical Layer, the Data Link Layer, the Routing Layer, the Transport Layer and the Application Layer. These layers correspond roughly to the layers as defined within the TCP/IP Suite. (The 802 standard has a different organizational structure for the layers and uses somewhat different names and numbering conventions.)\nAn important ideal in layered standards is the ideal of layer independence. A layered protocol suite specifies standard interfaces between layers such that, in theory, a device and protocol operating at one layer can coexist with any number of different protocols operating at higher or lower layers, so long as the standard interfaces between layers are followed.\nOptimizing The Handling of Data Units (Packets) within a Network Device\nWith an increasing number of network users needing increasing amounts of network bandwidth, there is a need for network devices to become efficient and sophisticated in their internal processing and routing of data. One mechanism known in the art for processing packets more efficiently is for network devices to identify flows in packets. A flow may be defined generally as a group of related packets that share some common characteristic which the network device may use to optimize the handling of the packets. For example, a network device may establish separate buffers for different flows so that very heavy traffic on one flow will not be allowed to cause congestion on another flow, or a network device may otherwise optimize its operation by grouping packets in flows.\nIn prior art systems, first the packets themselves are identified. The identification of packets is done based on the packet header. Each protocol has a header that uniquely identifies it. As described above, communication protocols frequently have encapsulated packets. That is, one protocol carried within another protocol. These are called protocol layers. In prior art systems, in order to group related packets together into flows, it is necessary to determine and identify multiple layers of protocols. The flow can be coarsely defined at a higher layer, requiring the identification of the outer protocol(s). A flow can also be defined at a finer layer of granularity by identifying multiple layers of encapsulated protocols.\nOnce a flow is identified, packets belonging to that flow can be isolated from other packets. The isolation of a flow permits the network device to process the flow independently from other packets. It also shields the flow from harmful effects caused by other flows.\nWithin the art it is becoming common to construct networking devices that perform layer three routing tasks using dedicated logic circuits. These devices are sometimes referred to as layer three switches. One task that these devices attempt to perform using fast dedicated logic is identifying and assigning packets to flows. In one such known system, for example, flows are identified by the system looking for packets flowing from a particular source address to a particular destination address. If a specified number of packets are received in a specified period of time having the same source and destination address, the system will establish a flow for those packets and may attach a flow identification field to those packets, allowing special handling of those packets within that device or by other networking devices in the network.\nThere is a trade-off, however, to be had between very fast layer three routing (or other processing of data units) done by dedicated hardware, and the flexibility that may be required in different networking operations. This flexibility is often achieved by using modifiable software to control data handling. Modern networks are often called on to carry a wide variety of network traffic under a wider variety of protocols, all over the same network hardware. Therefore, a networking device hard-wired to determine a flow based on a source and destination address at a IP or IPX layer may not be able to establish a different flow should it be desired based on an Ethernet address of a packet. The device also may not be able to be reconfigured to handle the inevitable changes to network protocols or new protocols that are developed.\nWhat is needed is a mechanism that would allow a high performance communication device to perform routing or sorting functions very quickly in hardware while allowing the operation of the routing function and flow assignment be specified and configured programmably so that it can be easily changed and can flexibly respond to different requirements of different network protocol layers."}
{"text": "As a way to combat global warming, introduction of new energy such as solar photovoltaic power generation and wind power generation has been promoted in recent years throughout the world. Since outputs of these power generations are affected by the weather, it is predicted that introduction on a large scale will cause problems with operation of power systems such as difficulty in maintaining frequencies and voltages. As a way to solve such problems, installation of large-capacity storage batteries for smoothing output variations, storing surplus power, and load leveling is expected.\nA redox flow battery is one of large-capacity storage batteries. In a redox flow battery, a positive electrode electrolyte and a negative electrode electrolyte are supplied to a battery cell having a membrane interposed between a positive electrode and a negative electrode, to charge and discharge the battery. An aqueous solution containing a metal ion having a valence which changes by oxidation-reduction is representatively used as the electrolytes. Representative redox flow batteries include an iron-chromium-based redox flow battery using an iron ion for a positive electrode and a chromium ion for a negative electrode, and a vanadium-based redox flow battery using a vanadium ion for both electrodes (e.g., Patent Document 1)."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a liquid ejection head for ejecting a liquid.\n2. Description of the Related Art\nA typical example of the liquid ejection head for ejecting a liquid includes an inkjet recording head applied to an inkjet recording system ejecting ink to a recording medium for recording. The inkjet recording head generally includes an ink flow path, an ejection energy generating unit disposed in a part of the flow path, and a fine ink ejection port for ejecting ink by energy generated therein.\nJapanese Patent Application Laid-Open No. 2006-168345 discloses a method for manufacturing a liquid ejection head applicable to the inkjet recording head. The method includes forming a flow path mold on a substrate having a plurality of ejection energy generating units; and applying thereto a covering resin layer made of a curable resin and for serving as a flow path wall member having a flow path wall. The method further includes curing a portion surrounding the mold serving as the flow path wall and containing an upper surface of the covering layer; laminating a polished silicon plate thereon; forming an ejection port in the plate; and then removing an uncured portion of the covering layer and the mold to form a space serving as the flow path.\nIn recent years, the recording apparatus has been needed to increase the image quality and the recording speed to higher levels. In order to meet the needs, the ejection ports and the flow path communicatively connected thereto are required to be arranged at high density and the volume of each droplet to be ejected is also required to be equalized at further higher levels.\nUnfortunately, the method disclosed in Japanese Patent Application Laid-Open No. 2006-168345 has a possibility that the covering layer may have a slightly uneven upper surface thereon due to existence of a partially remaining flow path mold in the entire substrate surface. If the silicon plate is arranged so as to conform to the uneven surface, the distance between the ejection energy generating unit and the ejection port may vary. If that happens, it is concerned that the variation of the distance causes a variation of the volume of a droplet ejected from each ejection port, which may affect an image to be recorded. Even if a high pressure is applied to laminate the silicon plate to the covering layer, it is difficult to planarize the covering layer sufficiently enough to remove the unevenness thereof because the upper surface of the covering layer is partially cured."}
{"text": "Provisioning a computing system can be a complicated and error-prone task. The computing system may need to be provisioned with many applications, each of which may need to have settings configured, data entered, and/or files altered. Provisioning multiple computing systems across an entire data center or even several data centers may be significantly more complicated and may involve the tedious and repetitive task of configuring the same enterprise-wide settings again and again. To reduce the complexity and rate of error, many provisioning systems now enable administrators to provision new systems from stored backup images of a fully configured system. Provisioning a new system from a backup image may allow administrators to quickly provision new systems both locally and remotely with a much reduced chance of configuration errors.\nUnfortunately, traditional systems for provisioning new systems from backup images may suffer from a variety of inefficiencies. Some traditional systems may aggregate read requests to the backup system in order to improve efficiency. However, sub-optimal read request batch sizes may lead to slow responses from the backup system, increasing the amount of time it takes to provision a new system from a backup and leading to user frustration. Accordingly, the instant disclosure identifies and addresses a need for additional and improved systems and methods for dynamically adjusting batch request sizes."}
{"text": "The invention relates generally to floor surface treatment apparatus, more particularly to such apparatus embodied, for example, in a floor scrubber having a rotary brush or brushes for scrubbing the floor or a rotary pad for waxing and/or polishing the floor.\nThe invention is in a category similar to that of the coassigned U.S. patent Ser. No. 09/934,146 of William R. Stuchlik filed Aug. 21, 2001, published May 2, 2002 as Publication No. U.S. 2002-0050014-A1, which is incorporated herein by reference. The invention has been developed because of the need becoming recognized for apparatus in which the floor surface treatment unit (e.g. brush, brushes, pad) may be differently positioned, as for variation of its pressure on the floor during operation, for transport of the apparatus, and for facilitating servicing of the floor surface treatment unit."}
{"text": "Low cost electronic communications may be realized through use of preexisting, open networks, in particular the Internet. A private network may utilize the existing Internet infrastructure to reduce the cost of implementing and maintaining hardware and software to establish networked communication for a group of private users.\nVirtual private networks (“VPN”) entail one approach to use of the Internet, or other publicly available network, as an alternative to expensive, dedicated communications networks. A VPN may utilize packet-switched communications in a software-defined, secure network that resides within a larger, publicly switched network. A telecommunications provider provides access to the public network for members of the VPN. Thus, the VPN shares the public network for communications traffic. In effect, the cost of building and maintaining the public network is shared by the many users of the network.\nVPNs are particularly cost effective, for example, for highly mobile workforces and smaller companies. A telecommunications company can provide the network, or preexisting networks, such as the telephone network or the Internet, can be utilized via the use of tunneling software to interface to private components of the network.\nSuch private networks inherently expose themselves to security risks. The Internet is an intentionally open, unsecured communication environment. It is designed to be available to the general public, businesses, government agencies and non-profit organizations. This openness leaves the network vulnerable to attacks, and those private networks that use the Internet for communications similarly expose themselves to attack via the Internet.\nIn contrast to the needs of the Internet, a private network must generally limit access. The private network must protect itself form security risks, as well as limit access to private resources as needed. Such restrictions help to maintain the integrity of data.\nPrivate networks, e.g. intranets or extranets, that are implemented with Internet-based interconnections typically use a number of approaches to protect the private components of the network from public access. Protection mechanisms include features such as firewalls, access lists, host and application layer security, and other tools to limit access via the Internet to intranet resources.\nInternet-based communications employ well-established, widely known communication protocols, resulting in well-known weaknesses. These weaknesses may be exploited for illicit access to private networks utilizing the Internet for some, or all, of their communications. While equipment and software vendors attempt to standardize their Internet targeted products, deficiencies in the standardization process create further weaknesses that arise due to differing implementations by various vendors.\nMoreover, increasing the size of a private network generally increases opportunities to exploit that network. A larger network not only presents more avenues for attack, but also presents greater difficulty in tracking access privileges, updating security procedures, and preserving synchronization between security procedures.\nNumerous techniques that exploit network deficiencies have evolved. These techniques include hijacking of a host address, spoofing an address and denial-of-service attacks. In the last of these, the perpetrator generally attempts to shut down a network resource, such as a host, by flooding the resource with messages.\nVarious systems have evolved to protect network communications. Commonly, encryption techniques are employed to hide the contents of network communications traffic. Some methods mask the real Internet Protocol (“IP”) address of source and destination hosts by “tunneling” through hardware gateways. Tunneling systems generally, however, can reveal true addresses between gateways. Existing systems also typically fail to guard against denial-of-service attacks.\nFurther, existing systems may fail to mask communication traffic patterns. Systems that provide some masking of traffic typically are unsuitable for packet-switched, networked environments.\nSome techniques protect against denial-of-service attacks by deploying redundant copies of critical data that reside on servers. Implementations typically employ either majority voting on fully replicated data servers or distributed encoded redundant data across physical servers. Synchronizing the data on the redundant servers is a complex task. Further, the IP addresses of the servers are fully exposed, as is the profile of data traffic. Unless the number of replicated servers is sufficiently large, their collective ability to withstand denial-of-service attacks is limited.\nMoreover, data assurance methods typically do not provide security, and may even decrease network security. For example, applying channel coding to message bits and blocks does not provide any data assurance during failure of a route or path. Neither do existing methods of data encryption and authentication provide data assurance when data packets are lost due to interception or jamming.\nTraditional methods of providing data security against eavesdropping (such as keyed encryption) grew out of point-to-point or single user communication channel models. Most communications now take place over networks and require improved methods of assurance and security."}
{"text": "a. Field of Invention\nThis invention pertains to a drum made of several components and having a unique cross-sectional shape, and more particularly to a drum having generally square or rectangular cross sectional shape with rounded corners and including a self supporting plastic liner surrounded by a corrugated fiber shell.\nb. Background of the Invention\nVarious goods, including liquids, dry or powdery materials, or semi-liquid goods are preferably shipped and stored in drums. Preferably these drums must be light weight so that they reduce the overall shipping weight and costs, and so that they can be manipulated easily when they are empty. Another criteria considered in making the drums is that they should have a shape which is space efficient so that a large number of drums can be stored efficiently in a small amount of space either during shipping or storage. Over the years drums having various shapes have been tried in the industry. For example cylindrical drums have been tried which are strong structurally but are not space-efficient because when several such drums are disposed side by side, there is a considerable amount of space left therebetween. Square or rectangular containers are much more space efficient since they can be disposed side-by side with virtually no dead space therebetween. However structurally they are not as strong as cylindrical drums, and furthermore such containers present sharp vertical edges which are particularly prone to crushing and even rupture.\nAdditionally, it was found that it is advantageous to make a drum made of two components: a plastic inner lining disposed in a fiber shell. However the shell must be usually made strong enough to support the inner lining and therefore it is usually bulky and must include various lining reinforcing means."}
{"text": "The present application relates generally semiconductor devices, and more specifically to low resistance conductive contacts and their methods of production.\nAdvances in integrated circuit (IC) materials and design have yielded generations of ICs where successive generations have smaller and more complex circuits. As ICs evolve, the functional density (i.e.; the number of interconnected devices per unit area) has generally increased and the critical dimension (i.e., the minimum feature size) has decreased. While dimensional scaling improves performance, increases production efficiency, and lowers costs, it has also increased the complexity of processing and manufacturing."}
{"text": "In the design, processing and implementation of embedded capacitors for power decoupling in electronic appliances and systems, different techniques have been proposed for the implementation of the embedded capacitors. In the implementation of embedded capacitors, limitations may arise in relation to the integration of the embedded capacitors with power distribution networks in terms of performance and space utilization. For example, blind/buried vias are typically required in existing techniques, resulting in process limitations and increased manufacturing complexity. Furthermore, in existing techniques, vias for high current and high-frequency decoupling, respectively, must typically be provided, resulting in increased number of interconnection and solder joints.\nOften, the different existing techniques for the implementation of embedded capacitors are designed and limited to a specific process technology during manufacturing, which can impose further limitations of existing techniques due to lack of adaptability for different manufacturing process technology."}
{"text": "Golf balls, whether of solid or wound construction, generally include a core and at least a cover or outer coating. The properties of a conventional solid ball may be modified by altering the typical single layer core and single cover layer construction to provide a ball having at least one mantle layer disposed between the cover and the core. The core may be solid or liquid-filled, and may be formed of a single layer or one or more layers. Covers, in addition to cores, may also be formed of one or more layers. These multi-layer cores and covers are sometimes known as “dual core” and “dual cover” golf balls, respectively. Additionally, many golf balls contain one or more intermediate layers that can be of solid construction or may be formed of a tensioned elastomeric winding, which are referred to as wound balls. One piece golf balls are even available. The difference in play characteristics resulting from these different types of constructions can be quite significant. The playing characteristics of multi-layer balls, such as spin and compression, can be tailored by varying the properties of one or more of these intermediate and/or cover layers.\nMoreover, color in a golf ball, being a dominant visual feature, is also capable of positively contributing to and enhancing a golfer's game by improving the player's ability to focus on the golf ball when swinging a club and striking the ball. It is desirable that a golfer's eye be drawn to the ball easily. By keeping an eye on the ball, the golfer is able to remain focused on the immediate task at hand of maintaining hand-eye coordination and producing great balance during swing with consistent spine angle in order for the club face to strike the golf ball with just the right force, depending on the chosen club and desired distance. In this way, a golf ball's unique visual appearance can influence and improve the golfer's physical performance substantially.\nMeanwhile, golf balls that are attractive and exude superior quality will also positively contribute to the psychological aspects of a golfer's game by boosting the golfer's confidence and morale, thereby motivating, inspiring and ultimately peaking performance on the green. Additionally, visually superior golf balls may be spotted and located more easily on the golf course, thereby reducing a player's stress level which naturally translates into improved scores. Accordingly, golf ball manufacturers desire to incorporate color in golf balls in order to beneficially impact and improve both the physical and emotional/psychological aspects of a golfer's game.\nToward this end, there remains a need for golf balls having superior overall color appearance to the human eye. The present invention addresses and solves this problem."}
{"text": "Carrier modulation techniques are used to transmit information signals from one location to another. Traditional signal modulation techniques include, for example, amplitude modulation (AM), frequency modulation (FM), phase modulation (PM). In addition, complex modulation techniques exist that incorporate aspects of AM, FM, and PM such as quadrature phase shift keying (QPSK), amplitude phase shift keying (APSK) and including quadrature amplitude modulation (QAM)."}
{"text": "1. Field of the Invention\nThis invention relates to building modules for the erection of walls, more particularly for freestanding, sound absorbing/reducing walls and retaining walls. These are particularly modular elements comprising at least two subelements, preferably in the form of beams arranged at an angle to one another. The modules may be regarded also as a building module kit from which at least one or more complete modules can be assembled. Accordingly, the term \"building module\" is intended to comprise unitary modules as well as multi-pan and more complex modules.\nIn many cases, one of the subelements or beams is constructed as a longitudinal beam extending substantially parallel to the wall plane and comprising at least two cross-beams arranged at an angle to one another. The cross beams extend between a pair of the longitudinal beams, in a direction transverse to the wall plane. Two of said cross-beams form support blocks with planar top and bottom sides which act as spacing means between the individual modules. Each module is designed to be held in place by gravity, its own weight, and the weight of a bulk filling material. Each module normally stacks upon another module to build up wall height, and is installed adjacent to another in an end-to-end arrangement to build wall length.\n2. Background Art\nThe typical retaining wall or freestanding wall is usually constructed with a plurality of flame-like elements comprising at least two subelements, preferably in the form of beams arranged at an angle to one another and connected in a form locking or material locking manner. U.S. Pat. No. 4,384,810, to Neumann, describes a locking beam that forms a three dimensional lattice in a construction system for plantable shoring walls. Said lattice comprises support blocks with planar top and bottom sides which are stacked one above the other. The blocks act as spacing means between individual planting levels. The structure includes a longitudinal component consisting of a base plate and a breast part. The longitudinal component and the locking beam extend parallel to the wall plane and always rest on two support blocks. The support blocks are spaced apart a distance and extend in a direction transverse to the wall plane. Each of the consecutive stacks of the support block pairs, together with the longitudinal component, the locking beam, and the earth within, form a construction section acting as a static slope shoring unit.\nU.S. Pat. No. 5,017,050, to Jaecklin describes building elements for supporting a grid wall with a bulk filling material comprising at least two subelements, preferably in the form of beams arranged at an angle to one another, and connected in a form locking or material locking manner.\nThere is at least one hole or recess in the cross beam which is open over a part of its circumferential contour and adapted to receive a longitudinal beam in order to establish a form locking connection between the beams. The beams are secured against separation and displacement from one another in a mounted state. The contour of the hole or recess overlaps the other side of the longitudinal beam so as to form an abutment. The recess is shaped so as to permit a partial lateral insertion of the longitudinal beam while establishing contact between the abutment and the longitudinal beam in a first rotational position. The longitudinal beam is then rotated downward through an arc until it is fully inserted into the recess. Once fully inserted, a form locking connection is established which cannot be reversed when the construction is under load.\nU.S. Pat. No. 5,181,351 to Jaecklin, describes building elements for supporting grid walls with a bulk material filling. The frame-like elements of the invention comprise at least two subelements, preferably in the form of beams arranged at an angle to one another, and connected in a form locking or material locking manner.\nOne of the subelements or beams is constructed as a longitudinal beam extending substantially parallel to the wall plane and comprises at least two profile legs arranged at an angle to one another. The first of these profile legs forms at least one bearing surface for the bulk filling material, while the second of these profile legs forms a retaining surface for the bulk filling material facing the inside space of the frame.\nSuch structural elements or structural systems have heretofore not been concerned with the need for structural features which allow for the adequate supply and retention of water necessary to sustain maximum plant growth without costly replanting. It is to this need and to other deficiencies of the prior art that this invention is directed."}
{"text": "1. Field of the Invention\nThe present invention relates to a constant current circuit and a reference voltage circuit using the same, and more particularly, to stabilizing operation of a constant current circuit.\n2. Description of the Related Art\nA conventional constant current circuit is described. FIG. 9 is a circuit diagram illustrating a conventional constant current circuit using the difference in K-value (drivability). The K-value is determined by K=W/L·(μCox/2), where W is the gate width, L is the gate length, μ is the mobility of carriers, and Cox is the gate oxide capacitance per unit area.\nThe conventional constant current circuit includes enhancement mode NMOS transistors 91 and 92 having different K-values, enhancement mode PMOS transistors 93 and 94, and a resistor 95.\nThe enhancement mode NMOS transistor 91 has a source terminal connected to a ground terminal 100 having a minimum potential, and a drain terminal and a gate terminal which are both connected to a gate terminal of the enhancement mode NMOS transistor 92 and a drain terminal of the enhancement mode PMOS transistor 93. The enhancement mode NMOS transistor 92 has a source terminal connected to the ground terminal 100 via the resistor 95, and a drain terminal connected to a gate terminal and a drain terminal of the enhancement mode PMOS transistor 94 and a gate terminal of the enhancement mode PMOS transistor 93. The enhancement mode PMOS transistors 93 and 94 each have a source terminal connected to a power supply terminal 101 having a maximum potential.\nNext, an operation of the conventional constant current circuit is described. The K-value of the enhancement mode NMOS transistor 91 is smaller than the K-value of the enhancement mode NMOS transistor 92. A voltage difference between a gate-source voltage of the enhancement mode NMOS transistor 91 and a gate-source voltage of the enhancement mode NMOS transistor 92 is generated across the resistor 95. A current flowing through the resistor 95 is mirrored by the enhancement mode PMOS transistors 93 and 94, thereby generating a bias current (see, for example, Japanese Patent Application Laid-open No. Hei 03-238513).\nHowever, the conventional constant current circuit has two operating points. One is a normal operating point at which the bias current flows. The other is an operating point at which the bias current becomes 0. When a potential at a connection point 291 becomes the maximum potential of the power supply terminal 101 and a potential at a connection point 290 becomes the minimum potential of the ground terminal 100, the constant current circuit is fixed at the operating point at which the bias current becomes 0, and thus fails to operate. The conventional constant current circuit has therefore a problem of needing a separate start-up circuit for start-up.\nIn addition, when the potential of the power supply terminal 101 increases and then the potential at the connection point 291 increases, the characteristics of the enhancement mode NMOS transistors 91 and 92 are changed by the channel length modulation effect of the enhancement mode NMOS transistor 92, with the result that the bias current fluctuates. In other words, the conventional constant current circuit has a problem of poor line regulation."}
{"text": "Personal audio devices, including wireless telephones, such as mobile/cellular telephones, cordless telephones, mp3 players, and other consumer audio devices, are in widespread use. Such personal audio devices may include circuitry for driving a pair of headphones or one or more speakers. Such circuitry often includes a speaker driver including a power amplifier for driving an audio output signal to headphones or speakers."}
{"text": "It is common for mass spectrometry instruments to have various regions operating at different pressures along the length of the instrument. For instance, the source region may be operated at relatively high or atmospheric pressure, whereas the analyser region may be operated under high vacuum conditions (i.e. very low pressures). Various regions, for instance regions used for transportation and/or manipulation of ions, may be provided between the source and analyser regions that are operated at intermediate pressures.\nFor example, WO 2013/171495 (MICROMASS) discloses a mass spectrometer in which ions generated at an ion source are passed in sequence through a first vacuum chamber via an ion transfer ion guide, into a second vacuum chamber containing a quadrupole mass filter, into a gas collision cell, into an ion mobility separation device and a second gas collision cell before finally arriving at the mass analyser, with each of these devices/regions being operated at various pressures. The operating pressure in each region is selected according to the purpose of that region.\nOther instruments having multiple pressures along the length of the instrument are disclosed e.g. in US 2013/0175442 (AGILENT), WO 00/16375 (VG ELEMENTAL), JP 2012-043672 (SHIMADZU) and WO 2014/080493 (SHIMADZU).\nIn some cases background gas may be allowed to flow freely between the different pressure regions. This gas may be used to carry ions in the direction of gas flow. For example, this may be exploited in the first vacuum stages of a mass spectrometer. In these instances the gas composition in sequential regions may be very similar.\nIn other cases, gases may be controlled such that they are prevented from entering another region by an outflow of gas from that region. For example, where ion mobility separation (“IMS”) devices are operated in the first vacuum stages of a mass spectrometer, ingress of source gas into the ion mobility separation device is prevented by excess pressure/flow of the ion mobility separation buffer gas. In this case, the ion mobility separation gas enters the previous or upstream region. This may also be the case where a high pressure ion mobility separation device is situated downstream of a low pressure quadrupole region.\nIt will be appreciated that in each of these cases there is an intermixing of gases between adjacent regions.\nIt is desired to provide improved methods of coupling gas-filled regions."}
{"text": "The present invention relates to a semiconductor variable capacitance element whose capacitance is controlled by an electric charge accumulated on a floating electrode which is covered with an insulator film provided in the surface of a semiconductor substrate so that the floating electrode is insulated from the outside.\nSemiconductor variable capacitance elements have heretofore been described in the literature (see Proceedings vol. 2 in 11th International Congress of Chronometry, '84, edited by The French Society of Microtechnology and Chronometry, p.9).\nFIG. 2 is a sectional view showing the structure of a conventional semiconductor variable capacitance element. On the surface of a semiconductor substrate 21 are provided a capacitance or terminal electrode 12 connected to an n-type diffused region 22, a floating electrode 23 having opposed major surface covered with an insulator film 25 so as to be insulated from the outside, and a variable or control electrode 13 for adjusting the amount of electric charge accumulated in the floating electrode 23.\nIn this semiconductor variable capacitance element, a depletion layer or capacitive region 24 formed in a surface region of the semiconductor substrate 21 which is located under the floating electrode 23 is controlled by a potential generated at the floating electrode 23 by the electric charge accumulated in the floating electrode 23. Since the floating and capacitance electrodes 23 and 12 are strongly capacitance-coupled to each other, the potential of the floating electrode 23 is varied in accordance with the voltage applied to the capacitance electrode 12.\nSince the capacitance electrode 12 of the conventional semiconductor variable capacitance element is connected directly to an external circuit, a bias voltage which is generated in the external circuit is applied directly to the capacitance electrode 12. For this reason, any variation in the voltage applied from the external circuit causes a change in the bias voltage, and this leads to variations in capacitance of the variable capacitance element. When the voltage applied from the external circuit is relatively high, a relatively high voltage is applied between the capacitance and floating electrodes 12 and 23, thus causing a minute tunnel current to flow into the floating electrode 23 through the insulator film 25. In consequence, the amount of electric charge accumulated on the floating electrode 23 is changed, and the capacitance value of the variable capacitance element is also gradually changed with time, disadvantageously."}
{"text": "1. Field of the Invention\nThe present invention relates to heat-accumulating microcapsule dispersion, and more particularly to a microcapsule dispersion including, in a stably dispersed state in heat transfer fluid medium, a number of heat-accumulating microcapsules each having a microcapsule shell accommodating therein organic compound functioning a heat accumulating material in association with a phase change thereof. The heat-accumulating microcapsule dispersion may be used as heat transfer medium employed in e.g. an air-conditioning system such as a local area air-heating system or an indoor air-cooling system.\n2. Description of the Related Art\nThere is known heat-accumulating microcapsule dispersion of the above-noted type including, in a stably dispersed state in water, a number of heat-accumulating microcapsules made of e.g. melamine resin having microcapsule shell accommodating therein heat-accumulating material such as tetradecane, paraffin wax or the like.\nFor manufacturing such heat-accumulating microcapsule dispersion as above, the heat-accumulating material and prepolymer of melamine resin are polymerized with each other while being dispersed and emulsified in water. As a result, there is obtained the dispersion in which heat-accumulating microcapsules each having a core formed mainly of the heat-accumulating material covered with a capsule outer layer microcapsule shell of the resin coating are dispersed in a stable manner in the water.\nThe heat-accumulating microcapsule dispersion of the above-noted type has a greater viscosity than e.g. water alone, because the dispersion includes the microcapsules dispersed therein. And, this viscosity tends to increase with use of the dispersion.\nOn the other hand, as for the heat-accumulating capacity, the dispersion has a higher capacity than the heat transfer fluid medium alone. Then, in order to obtain a certain fixed amount of heat transfer capacity, this is possible simply by circulating a smaller amount of dispersion in a circulating passage provided between a heat receiving end and a heat supplying end. That is, even if the diameter of the pipe constituting this passage is reduced, the same amount of heat transfer capacity may be obtained.\nAs shown in FIG. 2, the heat-accumulating microcapsule dispersion described above is responsible for transferring heat between a heat-receiving end heat exchanger and a heat-supplying end heat exchanger. Then, in order to allow the heat exchange to take place efficiently, one will encounter the problem of the heat transfer performance between the dispersion and the inner wall of the passage constituting the heat exchanging passage in which the dispersion runs.\nIn general, the heat transfer rate under such environment as above is a function of the Reynolds number which represents flow conditions of fluids. A smaller Reynolds number results in lower heat transfer rate.\nThe above will be described more specifically with reference to FIG. 6. In this FIG. 6, the horizontal axis represents the Reynolds number: Re of the dispersion flow, while the vertical axis represents the heat transfer rate: hi between the passage wall and the dispersion. An alternate long and short dashed line indicates the relationship between the Reynolds number: Re of water conventionally employed and the heat transfer rate: hi. A white round dot chain line indicates the relationship for the conventional dispersion including the conventional heat-accumulating microcapsules (the microcapsule has a volume average particle diameter smaller than 5 .mu.m and exhibits a particle distribution pattern as depicted in FIG. 4) dispersed in water.\nReferring further to the same figure, the mark: Reynolds (near 10000), denotes a typical Reynolds number obtained with a system operation using water alone as the operational medium. The further mark: Renew (near 1250), denotes a typical Reynolds number obtained with a system operation using dispersion including microcapsules dispersed therein in a stable manner.\nAs may be understood from this FIG. 6, in the case of the use of water alone or the conventional microcapsule, the heat transfer rate decreases with decrease in the Reynolds number. The heat transfer rate: hiold, of the case using water alone at Reold, is significantly higher than the heat transfer ratio: hinew of the case using the conventional microcapsule dispersed at Renew. In this respect, there is room for improvement.\nThat is to say, in case the heat-accumulating capsule is used as being stably dispersed in the heat transfer fluid medium in order to secure a certain heat transfer amount, it is desired to increase the heat transfer rate between this dispersion and the wall of the heat exchanging passage through which the dispersion is caused to flow.\nAccordingly, in view of the above-described drawbacks of the conventional art, a primary object of the present invention is to obtain heat-accumulating microcapsule dispersion which can provide a higher heat transfer rate between this dispersion and the member constituting the passage through which the dispersion is caused to flow."}
{"text": "The present invention relates to a device for controlling fuel injection quantity in diesel engines which are used as outboard marine engines.\nIn a typical diesel engine, the fuel injection device is controlled by a rotary control pinion and rack, and the quantity of injection of fuel from the fuel injection device is increased or decreased as the control pinion is rotated in one and the other directions.\nThis rotary control pinion, however, does not have a function of a safety device which, when a trouble has happened to take place in the diesel engine, automatically reduces the fuel injection quantity to finally stop the diesel engine.\nOn the other hand, in the case of diesel engine used as an outboard marine engine in which the crankshaft is disposed vertically, it is possible to lay a control rod for controlling the fuel injection quantity in parallel with the crankshaft and move the same in the vertical direction."}
{"text": "Variable valve activation (VVA) mechanisms for internal combustion engines are well known. It is known to lower the lift, or even to provide no lift at all, of one or more valves of a multiple-cylinder engine, during periods of light engine load. Such valve deactivation or valve lift switching can substantially improve fuel efficiency.\nA Roller Finger Follower (RFF), as a type of rocker arm, acts between a rotating eccentric camshaft lobe and a pivot point on the engine, such as a Hydraulic Lash Adjuster (HLA), to open and close an engine valve. Switchable RFFs may be a “deactivation” type or a “two-step” type. The term switchable deactivation RFF, as used herein, means the switchable RFF is capable of switching from a valve lift mode to a no lift mode. The term switchable two-step RFF, as used herein, means the switchable RFF is capable of switching from a first valve lift mode to a second and lesser valve lift mode that is greater than no lift. When the term “switchable RFF” is used herein, by itself, it includes both types.\nA typical switchable RFF includes an outer arm and an inner arm. The inner arm is movably connected to the outer arm. It can be switched by a locking member, from a coupled mode wherein the inner arm is immobilized relative to the outer arm, to a decoupled state wherein the inner arm can move relative to the outer arm. Typically, the outer arm of the switchable RFF is pivotally supported at a first end by the HLA. A second end of the outer arm operates against an associated engine valve for opening and closing the valve by the rotation of an associated eccentric cam lobe acting on an inner arm contact surface which may be a roller. The inner arm is connected to the outer arm for pivotal movement about the outer arm's second end with the contact surface of the inner arm disposed between the first and second ends of the outer arm. Typically, the locking member includes a locking pin disposed in a bore in the first end of the outer arm, the locking pin being selectively moved to engage the inner arm to thereby couple the inner arm to the outer arm when engaged, and decouple the inner arm from the outer arm when disengaged.\nIn a switchable two-step RFF, the outer arm typically supports a pair of rollers carried by a shaft. The rollers are positioned to be engaged by associated low-lift eccentric cam lobes that cause the outer arm to pivot about the HLA, thereby actuating an associated engine valve to a low-lift. The inner arm, in turn, is positioned to engage an associated high-lift eccentric cam lobe sandwiched between the aforementioned low-lift lobes. The switchable two-step RFF is then selectively switched between a coupled and a decoupled mode by the locking member. In the coupled mode, with the inner arm locked to the outer arm, the rotational movement of the central high-lift lobe is transferred from the inner arm, through the outer arm to cause pivotal movement of the RFF about the HLA, which, in turn, opens the associated valve to a high-lift. In the decoupled mode, the inner arm is no longer locked to the outer arm and is permitted to move relative to the outer arm against a lost motion spring that biases the inner arm away from the outer arm. In turn, the rollers of the outer arm engage their associated low-lift lobes. The rotational movement of the low-lift lobes is transferred directly through the outer arm, and the associated valve is reciprocated by the outer arm to a low-lift.\nA switchable deactivation RFF typically includes an outer arm and an inner arm. The inner arm supports a roller carried by a shaft. The roller is engaged by an eccentric lifting cam lobe for actuating an associated engine valve. Like the switchable two-step RFF, the switchable deactivation RFF is selectively switched between a coupled and a decoupled mode by a movable locking member. In the coupled mode the inner arm of the switchable deactivation RFF is locked to the outer arm and the rotational movement of the associated lifting cam lobe is transferred from the inner arm, through the outer arm to cause pivotal movement of the RFF about the HLA, which, in turn, opens the associated valve to a prescribed lift. In the decoupled mode, the inner arm becomes unlocked from the outer arm and is permitted to pivot relative to the outer arm against a lost motion spring. In the decoupled mode, the rotational movement of the lifting cam lobe is absorbed by the inner arm in lost motion and is not transferred to the outer arm. Thus, the associated valve remains closed when the switchable deactivation RFF is in its decoupled mode.\nIn a first switchable deactivation RFF design, the inner arm makes contact with an associated cam lobe while the outer arm does not. The lost motion spring biases the inner arm away from the outer arm and, with the outer arm supported by the HLA, serves to load the inner arm against its associated cam lobe in the decoupled mode. In a switchable deactivation RFF having a lost motion spring with an effective force exerted on the HLA that is higher than the opposing force of an associated HLA spring, the opposing forces must be properly managed to prevent reactive pump-down of the HLA induced by the force of the lost motion spring. For this purpose, an expansion travel limiter is incorporated in the switchable deactivation RFF to limit the movement of the inner arm relative to the outer arm. Thus, when the switchable deactivation RFF is in its decoupled mode, and the inner arm of the RFF follows the cam lobe, the lost motion spring will push the outer arm until the expansion travel limiter is engaged. At that point, further movement of the outer arm relative to the inner arm ceases, HLA pump-down is prevented and HLA leak-down recovery is initiated. Moreover, at that point, since the effective preload force of the lost motion spring is greater than the expansion force of the HLA, pump-up of the HLA is prevented. Note also that, when the inner arm of the RFF follows the base circle of the associated cam lobe, the expansion travel limiter further serves to set a clearance gap or mechanical lash between the locking pin and the inner arm to assure proper alignment of the locking pin with the inner arm when the RFF switches between its decoupled and coupled modes and to define the total mechanical lash in the valve train.\nIn an alternate switchable deactivation RFF design, in order to increase its resistance to HLA pump-up beyond that provided by the installed load of the lost motion spring, null pads may be added on the outer arm of the switchable deactivation HLA for contacting zero-lift, constant radius null lobes disposed on either side of the associated lifting lobe. In this design, when the inner arm contacts the expansion travel limiter, the inner arm of the RFF is prevented from contacting the base circle of its associated cam lobe by the null pads first contacting with the zero-lift null lobes. Since the inner arm is held away from the base circle of the cam by the expansion travel limiter, the force of the lost motion spring cannot pump-down the HLA. By the null pads making contact with the zero-lift null pads, pump-up of the HLA is prevented by the opposing installed load of the associated valve closing spring. The expansion travel limiter establishes the mechanical lash between the locking pin and the inner arm; as well as the clearance (lash) between the inner arm and the base circle of the cam. The pin lash plus the cam lash establishes the total mechanical lash of the valve train.\nVarious lost motion expansion limiters used in switchable RFFs are known in the art. For example, in U.S. Pat. No. 6,532,920, a switchable two-step RFF is shown wherein the roller shaft of the outer arm contacts a throughbore in the inner arm to limit inner arm travel. As shown in U.S. Pat. Nos. 5,544,626, 5,653,198 and 6,314,928, bumper pads or projections formed at the lower end of the inner arm are used to limit inner arm travel of the switchable RFFs. The disadvantage of these devices in the prior art is that the stop position cannot be precisely controlled resulting in sometimes too small or too large of a mechanical lash between the locking pin and the inner arm or, in the case of a switchable deactivation RFF with null pads, resulting in a clearance between the inner arm and base circle of the associated cam lobe that is too too small, or even non-existing. A mechanical lash that is too small may result in the locking pin being unable to reliably engage the inner arm. A lash that is too large may permit excess pump-down of the HLA thereby delaying the opening point, decreasing the lift and advancing the closing point of the associated valve in the coupled mode which is known to contribute to engine roughness at idle and/or emission problems. A cam clearance that is too small (in the case of a switchable deactivation RFF with null pads), between the inner arm and its base circle, increases total lash when the inner arm is allowed to contact the cam base circle and may similarly affect the opening, closing and lift characteristics of its associated valve.\nWhat is needed in the art is a device that precisely limits the amount of upward pivotable movement of the inner arm relative to the outer arm caused by the force of the lost motion spring.\nIt is a principal object of the present invention to provide an inner arm stop to precisely position the inner arm relative to the outer arm thereby controlling mechanical lash and HLA pump-down."}
{"text": "1. Field of the Invention\nThis invention relates to a process for the preparation of micro inorganic foamed granules.\n2. Description of Prior Arts\nInorganic foamed granules are lightweight, and have been accordingly used widely as lightweight aggregates for construction of buildings, extenders for fertilizer, detergent or paint, absorbents, soil improvers, etc. Particularly, the micro inorganic foamed granules having a mean particle size of not larger than 500 .mu.m are employed as fillers for imparting heat insulating properties to buildings and adjustors for adjusting specific gravity of cement used for an oil well or geothermal well, in addition to the above usage.\nParticularly, the use as adjustor for adjusting specific gravity of cement used for an oil well or geothermal well is one of the quite important uses.\nIt is required that the micro inorganic foamed granules for employing as specific gravity adjustor should have apparent specific gravity in the range of 0.6 to 0.9 and particle size of not larger than 500 .mu.m (preferably not larger than 200 .mu.m). Further, high strength is required, and the micro inorganic foamed granules having hydrostatic pressure strength of not lower than 300 kgf/cm.sup.2 at 70% of indestructibility are generally employed for the above purpose.\nThe known micro inorganic foamed granules are classified into three groups, that is, granules manufactured from natural volcanic glassy materials, those manufactured from artificial glassy materials and those contained in coal ash.\nThe micro inorganic foamed granules of natural volcanic glassy materials are produced by the steps of pulverizing materials such as obsidian, pearlite and volcanic ash to make the particle size not larger than 100 .mu.m, removing a portion of fine powder, adding thereto additives such as SiO.sub.2 powder and Al.sub.2 O.sub.3 powder in an appropriate amount and firing the mixture at a temperature not lower than 1,000.degree. C.\nThe micro inorganic foamed granules of natural volcanic glassy materials prepared in the above method generally has a low apparent specific gravity and can be manufactured at relatively low cost. However, hydrostatic pressure strength thereof is extremely low. Therefore, the micro inorganic foamed granules of natural volcanic glassy materials cannot be employed as adjustor for adjusting the specific gravity of cement used for an oil well or geothermal well. The foamed granules are used exclusively as heat insulating materials or lightweight aggregates for construction of buildings, which do not require mechanical strength.\nThe micro inorganic foamed granules of artificaial glassy materials are produced by steps of adding a foaming compound such as carbon to artificially prepared glassy materials, fusing the mixture and spraying.\nThe above micro inorganic foamed granules are easily manufactured so as to have the particle size of not larger than 100 .mu.m. The apparent specific gravity and hydrostatic pressure strength of the obtained micro inorganic foamed granules are satisfactory enough to employ as adjustors for adjusting specific gravity of cement used for an oil well or geothermal well. However, manufacturing cost thereof is high.\nInorganic foamed granules contained in coal ash in a very small amount are known. The inorganic foamed granules show favorable properties such as 0.6 to 0.7 of apparent specific gravity and 400 kgf/cm.sup.2 of hydrostatic pressure strength (measured at 70% of indestructibility) for employment as adjustor for adjusting specific gravity of cement used for an oil well or geothermal well. However, the foamed granules are generally contained in coal ash in the amount of not more than 0.5 wt.%. For this reason, the foamed granules are not collected industrially in Japan, which makes the granules very expensive."}
{"text": "1. Field of the Invention\nThe present invention relates an image processing apparatus and an image processing method for combining images.\n2. Description of the Related Art\nConventionally, an image processing apparatus that combines a plurality of images has been available to generate a single image. For combining a plurality of images, such an image processing apparatus discards attribute information attached to the original images or selects the attribute information of one of the images.\nHowever, when the attribute information is processed as described above, optimum attribute information is not supplied to a combined image. This makes it impossible to switch processing for optimum image processing in accordance with attribute information."}
{"text": "Modern data processing systems, such as general purpose computer systems, allow the users of such systems to create a variety of different types of data files. For example, a typical user of a data processing system may create text files with a word processing program such as Microsoft Word or may create an image file with an image processing program such as Adobe's PhotoShop. Numerous other types of files are capable of being created or modified, edited, and otherwise used by one or more users for a typical data processing system. The large number of the different types of files that can be created or modified can present a challenge to a typical user who is seeking to find a particular file which has been created.\nModern data processing systems often include a file management system which allows a user to place files in various directories or subdirectories (e.g. folders) and allows a user to give the file a name. Further, these file management systems often allow a user to find a file by searching for the file's name, or the date of creation, or the date of modification, or the type of file. An example of such a file management system is the Finder program which operates on Macintosh computers from Apple Computer, Inc. of Cupertino, Calif. Another example of a file management system program is the Windows Explorer program which operates on the Windows operating system from Microsoft Corporation of Redmond, Wash. Both the Finder program and the Windows Explorer program include a find command which allows a user to search for files by various criteria including a file name or a date of creation or a date of modification or the type of file. However, this search capability searches through information which is the same for each file, regardless of the type of file. Thus, for example, the searchable data for a Microsoft Word file is the same as the searchable data for an Adobe PhotoShop file, and this data typically includes the file name, the type of file, the date of creation, the date of last modification, the size of the file and certain other parameters which may be maintained for the file by the file management system.\nCertain presently existing application programs allow a user to maintain data about a particular file. This data about a particular file may be considered metadata because it is data about other data. This metadata for a particular file may include information about the author of a file, a summary of the document, and various other types of information. A program such as Microsoft Word may automatically create some of this data when a user creates a file and the user may add additional data or edit the data by selecting the “property sheet” from a menu selection in Microsoft Word. The property sheets in Microsoft Word allow a user to create metadata for a particular file or document. However, in existing systems, a user is not able to search for metadata across a variety of different applications using one search request from the user. Furthermore, existing systems can perform one search for data files, but this search does not also include searching through metadata for those files."}
{"text": "The present invention relates to a reader-printer having a device for rotating an image of microfilm on a projection surface.\nThe frames on a microfilm are in a vertical or horizontal position relative to the film, while information is recorded in the frame in a vertical or horizontal position relative to the frame. If images on the microfilm are projected with the film in a definite position, it is likely that the image will be projected in an inappropriate orientation. For example, information projected as oriented horizontally in the reader mode is almost impossible to read. Further if the orientation in which the frame is projected is not in coincidence with the orientation of the copy paper in the printer mode, an incomplete copy will be obtained. These objections can be eliminated by rotating the microfilm, but the device for holding the film must then be rotated. This requires a complex arrangement.\nU.S. Pat. No. 3,907,418 discloses an arrangement wherein the film image to be projected is rotated on a plane of projection by a prism. The prism employed is a trapezoidal prism, i.e., Dove prism, and is disposed at the light emergent side of a projection lens rotatably about the optical axis of the lens. The rotation of the prism rotates the image through twice the angle of rotation thereof. The arragement of U.S. Pat. No. 3,907,418 is based on the assumption that although the frames are directed vertically or horizontally relative to the film, information is recorded in the frames in an upright position relative to the film, and is adapted to selectively determine whether the prism in an orientation in the reader mode is to be rotated through 45 degrees in operative relation with the change of mode to the printer mode, such that the prism is rotated only when the direction of the frame is not in coincidence with the orientation of the paper in printer mode.\nHowever, the disclosed technique still requires the user's manipulation and has yet to be improved in the ease of use.\nThe technique has another problem in that the prism wherein the incidence optical axis coincides with the emergence optical axis is large-sized, whereas use of a small prism for compacting the apparatus involves deviation between the incidence optical axis and the emergence optical axis to result in the deviation of the projected image on copy paper."}
{"text": "Generally, in the pump driving mechanism, a pure reciprocating motion must be given to a piston to prevent water leakage and a contrivance to achieve that motion, such as a crank connection is required.\nTherefore, in such conventional pump driving mechanisms as the one which has been adopted in a water feeding device 70 of, for example, FIG. 15, a driving force of an output shaft 72 of an electric motor 71 for a pump is transmitted to a driving gear 74 through a gear 73. An eccentric rotary motion of an eccentric pin 75 provided to project on the driving gear 74 is converted to a reciprocating motion through a slider 76 and piston joint 77, and a piston 78 is correctly driven to reciprocate within a cylinder 79. In such mechanism, however, a large noise has been generated when it is driven and the energy loss has been large for the following reasons. That is, there are seven places in which the respective elements slide in contact with each other, namely, between the eccentric pin 75 and the slider 76, between the slider 76 and the piston joint 77, between the piston joint 77 and the driving gear 74, between the slider 76 and the driving gear 74, between a piston rod 80 and a shaft bearing 81, between the gear 73 and the driving gear 74 and between the driving gear 74 and a reduction gear shaft 82 and the noise has been large."}
{"text": "1. Field of the Invention\nThe present invention generally relates to methods and systems for generating simulated images from design information.\n2. Description of the Related Art\nThe following description and examples are not admitted to be prior art by virtue of their inclusion in this section.\nFabricating semiconductor devices such as logic and memory devices typically includes processing a substrate such as a semiconductor wafer using a large number of semiconductor fabrication processes to form various features and multiple levels of the semiconductor devices. For example, lithography is a semiconductor fabrication process that involves transferring a pattern from a reticle to a resist arranged on a semiconductor wafer. Additional examples of semiconductor fabrication processes include, but are not limited to, chemical-mechanical polishing (CMP), etch, deposition, and ion implantation. Multiple semiconductor devices may be fabricated in an arrangement on a single semiconductor wafer and then separated into individual semiconductor devices.\nAs design rules shrink, the design that is formed on a specimen such as reticles and wafers, even when formed using an optimally performing process, can look much different from the actual design. For example, due to the inherent limitations of the physical processes involved in forming a design on a physical specimen, features in the design formed on the physical specimen typically have somewhat different characteristics than the design such as different shapes (e.g., due to corner rounding and other proximity effects) and can have somewhat different dimensions (e.g., due to proximity effects) even when the best possible version of the design has been formed on the specimen.\nSometimes, it is not possible to know how the design will appear on the specimen and in images of the specimen, on which the design information has been formed, generated by tools such as inspection tools, defect review tools, metrology tools and the like. However, it is often desirable to know how the design will appear on the specimen and in images generated by such tools for a number of reasons. One reason is to make sure that the design will be formed on the specimen in an acceptable manner. Another reason is to provide a reference for the design, which illustrates how the design is meant to be formed on the specimen, that can be used for one or more functions performed for the specimen. For example, in general, a reference is needed for defect detection so that any differences between the design formed on the specimen and the reference can be detected and identified as defects or potential defects.\nMuch work has therefore been done to develop various methods and systems that can simulate how the design will be formed on the specimen and will appear in images of the specimen. There are several currently used methods for generating such simulated images. For example, one currently used method is forward electromagnetic (EM) modeling. In this method, design information such as computer aided design (CAD) and material information for a specimen is used as input, and the method simulates the physical interaction between light and material and generates a near field or far field image. Another example of a currently used method is rule-based approximation. Given some practical constraints, this method handcrafts a collection of rules that approximate transformation from CAD to observed image.\nThere are, however, several disadvantages to the currently used methods for generating simulated images. For example, forward EM simulation is extremely computationally intensive, making it infeasible for use in real production, e.g., the simulated images cannot be generated fast enough using forward EM simulation for the production use case. This is true for mask inspection, wafer inspection, and wafer metrology applications. In addition, forward EM simulation cannot simulate CMP and etching effects captured by electron beam and optical images. In another example, rule-based approximation tends to result in a large collection of very complex rules, making it generally inapplicable and causing poor prediction in reality.\nAccordingly, it would be advantageous to develop systems and methods for generating simulated images from design information that do not have one or more of the disadvantages described above."}
{"text": "Wireless optical input devices are well known in the art. A typical optical input device uses a light source to illuminate a portion of a navigation surface and a sensor to generate digital images from the light that is reflected off the navigation surface. In general, an optical input device tracks the relative movement between a navigation surface and the input device, multiple frames of digital image data of the illuminated navigation surface are captured by the sensor and successive digital images are subsequently processed and further translated as a cursor movement on the input device.\nAs the wireless input devices are becoming more popular, power saving has become one of the critical considerations in the design of an energy efficient wireless input device. It is desired to reduce the power consumption of a wireless input device during operation in order to increase the battery life of the device. Additionally, prolonging the battery life may reduce the adverse impact on the environment associated with the disposal of used batteries.\nOne of the methods that is commonly adopted to improve the energy efficiency of a wireless input device is to put the input device into a rest mode when it is inactive or stays idle for a predetermined period of time. When an input device is in the active mode, the input device operates at full operational capacity. On the other hand, when the device is in a rest mode, portions of the device circuitry are normally rendered inactive, thus, reducing the power consumption. However, in a conventional design, even when the input device is in the rest mode, a wake up detector on the input device may still be in operation, whereby the wake up detector is configured to trigger the input device to periodically check the sensor in order to monitor for movement and immediately wake up the input device if movement is detected. Hence, the input device is periodically being switched to active mode even though it is in an idle state. As a result, much of the supposedly inactivated device circuitry is being activated during the wake up detection operation. Therefore, the input device is actually consuming considerable battery power even though it is in a rest mode while inactive.\nIn general, during a typical wake up detection operation, the whole image data on the sensor array may be processed by the input device during the motion detection. However, since the computed motion data would not be outputted by the sensor to facilitate any navigation operation during a wake up detection operation. Therefore, an accurate motion data computation is actually not critical and may be unnecessary during the wake up detection operation. Accordingly, there is no need to process the whole image data on the sensor array during waking up detection operation and consuming battery power unnecessarily.\nThroughout the description and figures, similar reference numbers may be used to identify similar elements."}
{"text": "Laser diode arrays and arrays of laser diode bars are commercially available. Such arrays are fabricated typically by mounting individual laser diodes onto the top surface of a heat sink such that the emitting facet of the laser directs coherent light typically through a microlens positioned to collimate that light.\nWhen individual laser diodes are stacked, the heat sinks are separated by a layer of electrically insulating material having a thickness equal to that of the laser diode and the insulating layer as well as the top of the laser diode are covered by a layer of electrically conducting material to which the adjacent heat sink is connected. The individual laser diodes (components) are aligned with one another in the stack such that all the emitting facets are in a single plane which is the same plane occupied by the end faces of the heat sinks.\nIn such arrangements, adjacent components are placed as close together as possible. But the minimum spacing, q, between adjacent components (i.e. the laser diodes or the light beams emitted thereby) is determined by the thickness of the components and the heat dissipating capacity of the stack. The best available separation between adjacent components in the stack (or array), whether the components are individual laser diodes or laser diode bars, is 1.2 mm which is capable of producing a power density of 200 watts/sq cm (cw). A characteristic of prior art stacks is that q=p."}
{"text": "Imidazolinone compounds, for instance, those described in U.S. Pat. Nos. 4,188,487; 4,798,619 and 5,334,576, are highly potent, broad spectrum, environmentally benign, herbicidal agents. In general, the herbicidal activity of the R-isomer is better than that of the racemic imidazolinone compound. A process to prepare chiral imidazolinones via the resolved optically active 2-amino-2,3-dimethylbutyramide enantiomers is described in U.S. Pat. No. 4,683,324. Said aminoamide enantiomers are prepared via the hydrolysis of their chiral 2-amino-2,4-dimethylbutyronitrile precursors and are difficult to isolate.\nTherefore, it is an object of this invention to provide a stereospecific process to prepare chiral nicotinic, quinolinic or benzoic acid imidazolinone herbicidal agents directly from (R)2-amino-2,3-dimethylbutyronitrile without loss of optical purity and without the prior formation of (R)2-amino-2,3-dimethylbutyramide.\nThe present invention provides a stereospecific process to prepare a chiral compound of formula I \nwherein\nX is N or CH; and\nY and Z are each independently H, C1-C4alkyl optionally substituted with one C1-C4alkoxy group or Y and Z may be taken together to form a group xe2x80x94CHxe2x95x90CHxe2x80x94CHxe2x95x90CHxe2x80x94\nwhich process comprises the following steps:\na) reacting a compound of formula II \nwherein X, Y and Z are as described hereinabove with at least one molar equivalent of (R)2-amino-2,3-dimethylbutyronitrile in the presence of a non-polar, essentially water-free solvent, optionally in the presence of a tertiary amine, to form a first reaction mixture;\nb) hydrolyzing said reaction mixture in aqueous acid to form an acidic second reaction mixture;\nc) reacting said second reaction mixture with an excess of an aqueous base at a temperature of about 20xc2x0-85xc2x0 to form a basic third reaction mixture;\nd) separating said third reaction mixture to obtain an aqueous phase; and\ne) acidifying said aqueous phase to obtain the desired chiral formula I imidazolinone compound.\nChiral imidazolinone compounds having the R configuration demonstrate about a 2-fold increase in herbicidal activity over the corresponding racemic mixture. Heretofore, (R)imidazolinone compounds were prepared from (R)2-amino-2,3-dimethylbutyramide due to the instability of the (R)2-amino-2,3-dimethylbutyronitrile compound. However, isolation of said (R)aminoamide is difficult. Surprisingly, it has now been found that chiral imidazolinone herbicides may be prepared directly from (R)2-amino-2,3-dimethylbutyronitrile in the presence of a non-polar essentially water-free solvent with substantially complete retention of enantiomeric purity from the (R)aminonitrile starting material to the final chiral imidazolinone herbicidal product. Advantageously, the process of the invention eliminates the need for the prior formation of (R)2amino-2,3-dimethylbutyramide.\nIn accordance with the process of the invention, a formula II anhydride is reacted with at least one molar equivalent of (R)2-amino-2,3-dimethylbutyronitrile in the presence of a non-polar, essentially water-free solvent, optionally in the presence of a tertiary amine, to form a first reaction mixture; said reaction mixture is hydrolyzed with aqueous acid to form an acidic second reaction mixture; said second reaction mixture is treated with an excess of a base at a temperature of about 20xc2x0-90xc2x0 C. to form a basic third reaction mixture; said third reaction mixture is separated to obtain an aqueous phase; and the aqueous phase is acidified to obtain the desired chiral formula I imidazolinone herbicide. The process is illustrated in flow diagram I. In the specification and claims, an asterisk designates the assymetric carbon upon which the (R) configuration is conferred. \nSolvents suitable for use in the process of the invention are non-polar essentially water-free solvents such as aromatic hydrocarbons (e.g. toluene, benzene, xylene, naphthalene and the like, preferably toluene), halogenated aromatic hdrocarbons (e.g. chlorobenzene, dichlorobenzenes and the like), hydrocarbons (e.g. pentanes, hexanes and the like), halogenated hydrocarbons (e.g. chloroform, methylene chloride, dichlorethane, and the like, esters (e.g. ethyl acetate, methyl propionate and the like), ethers (e.g. diethyl ether, tetrahydrofuran, dioxane and the like) or any of the conventional, preferably water immiscible, organic non-polar solvents.\nPreferred non-polar solvents suitable for the process of the invention are aromatic hydrocarbons, particularly toluene.\nTertiary amines suitable for use in the first step of the inventive process are pyridine, 4-cyanopyridine, 4-picoline, 2-picoline, mixed picolines, tri(C1-C4)alkylamine, quinoline or any of the conventional organic tertiary amines, preferably 4-picoline. The amine may be present in amounts ranging from catalytic to excess amounts such as 10 mole % to 4.0 molar excess.\nAcids suitable for use in the process of the invention include strong mineral acids such as HCl or H2SO4, preferably H2SO4.\nBases suitable for use in the inventive process include alkali metal hydroxides or alkoxides, preferably hydroxides such as NaOH or KOH, preferably NaOH. These may be present at about 2 to 20 molar-equivalents, preferably about 2 to 8 molar-equivalents.\nIt is also intended that the process of the invention embraces the use of (S)-2,3-dimethylbutyronitrile to prepare the corresponding (S)-imidazolinone herbicidal product.\nIn actual practice, a mixture of the formula II anhydride in a non-polar, essentially water-free solvent, preferably an aromatic hydrocarbon, more preferably toluene, is treated with a 10% to 60% solution of (R)2-amino-2,3-dimethylbutyronitrile in a non-polar, essentially water-free solvent, preferably an aromatic hydrocarbon, more preferably toluene, optionally in the presence of 10 mol % to 4.0 molar excess, preferably about 10 mol % to 1.0 molar equivalent of 4-picoline, at a temperature of about 5xc2x0 to 45xc2x0 C., preferably about 5xc2x0 C. to 30xc2x0 C., to form a first reaction mixture; said mixture is treated with a strong mineral acid, preferably H2SO4, and water at temperatures of about 5xc2x0 to 80xc2x0 C., preferably about 20xc2x0 C. to 60xc2x0 C. to form an acidic second reaction mixture; said second mixture is treated with an excess, (about 2 to 20 moles per mole of acid used in the previous hydrolysis step) of an aqueous alkali metal hydroxide or alkoxide, preferably an alkali metal hydroxide, more preferably NaOH or KOH, of 10% or greater concentration on a weight basis at temperatures of about 15xc2x0 C. to 90xc2x0 C., preferably 20xc2x0 C. to 85xc2x0 C., to form a basic third reaction mixture; said third reaction mixture is separated to obtain an aqueous phase; and said aqueous phase is acidified with a strong mineral acid such as HCl, HBr or H2SO4, preferably H2SO4, to a pH of about 2 to 4 to obtain the desired chiral formula I imidazolinone product. The product may be isolated using conventional procedures such as filtration, extraction with a suitable solvent, chromatographic separation and the like, preferably filtration or extraction.\nIn order to facilitate a further understanding of the invention, the following examples are presented primarily for the purpose of illustrating certain more specific details thereof. The invention is not to be deemed limited thereby except as defined in the claims. Unless otherwise noted, all parts are by weight. NMR designates nuclear magnetic resonance. HPLC designates high performance liquid chromatography."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus which includes a transfer member for transferring toner on an image bearing member and presses the transfer member toward the image bearing member to form a transfer nip.\n2. Description of the Related Art\nConventionally, as a color image forming apparatus such as a printer and copier, there has been disclosed a technique in which toner images formed on respective photosensitive drums are sequentially superimposed on an intermediate transfer belt and the toner images superimposed on an intermediate transfer belt are collectively transferred onto a recording material to form a color image. In addition, there has been disclosed a technique in which toner images formed on respective photosensitive drums are sequentially superimposed onto the recording material on a conveying transfer belt using the conveying transfer belt which conveys the recording material while the recording material is electro-statically adhered, to form a color image.\nA transfer member transferring the toner images on the photosensitive drum is arranged at a position where the respective photosensitive drum faces toward the intermediate transfer belt or the conveying transfer belt. The transfer member is pressed against the photosensitive drum by a press member such as a spring to form a transfer nip. A transfer roller rotating to follow the intermediate transfer belt or the conveying transfer belt is generally used as the transfer member.\nHowever, the transfer roller is pressed by the press member at both axial ends of the transfer roller to produce a deflection recessed in a bow shape at the axial central portion. This causes the nip pressure of the transfer nip to be lower at its both axial ends than that at its central portion, making the nip pressure uneven particularly in an image forming region (an image bearing region where the toner images are borne), which may cause a failure in transfer. In order to cope with the problem, there is known a technique in which the outer diameter of the roller is made larger at its central portion than that at its both axial ends to uniform the nip pressure of the transfer nip even if the axis is deflected.\nHowever, the transfer roller is provided with an elastic member around its axis. An elastic force of the elastic member may vary the transfer nip secular change with time. The transfer roller is originally manufactured such that the outer diameter thereof is axially varied, complicating the production process."}
{"text": "1. Field of the Invention\nThe present invention generally relates to devices and apparatuses used to train athletes to hit a moving ball, and more particularly relates to pitching machines utilized by a coach to fire a ball towards a batter, with the pitching machine utilizing elastic bands or straps for propulsion of the ball.\n2. Background Information\nThe training of athletes, namely baseball and softball players, to properly hit baseballs/softballs is frequently done. Typically, a batter will spend time within a “batting cage” receiving balls pitched by a commercial pitching machine. Alternatively, portable pitching machines are known in the prior art and are utilized by coaches to pitch balls to batters.\nThere are many disadvantages in the prior art. A first disadvantage is the fact that balls are often fired along a generally constant flight path through use of a mechanized pitching machine that can be aimed differently and the motor on the machine can be sped up and/or slowed down. In addition, it is generally difficult to deliver to a batter balls pitched at a great variety of pitching speeds and pitching locations in real time.\nA second disadvantage to the utilization of commercial pitching machine equipment is cost. These commercial pitching machines typically utilize actual baseballs or softballs as the propelled ball. Such units are typically expensive to purchase, own, and operate. Additionally, such units typically require some type of netting to contain the pitched and hit balls.\nAn additional concern is that mechanical pitching machines are typically powered through the use of an electrical motor, thereby requiring the utilization of the unit in an area that has convenient access to an electrical power supply.\nWhat is needed is a less expensive, easy to use, adjustable, variable speed, non-motorized, portable pitching machine that is self-contained, allows an operator to shoot balls toward a ball player by using an elastic sling type pouch, is used for the purpose of hitting or fielding practice, and allows the operator to comfortably position himself or herself on top of the apparatus and behind the projection sling pouch for convenient and accurate operation. The apparatus allows the operator to easily adjust the speed and location of each pitch by simply adjusting the point of release of the ball pouch. The apparatus preferably utilizing perforated plastic balls. Embodiments of the present invention satisfy this need.\nAdditional objects, advantages and novel features of the invention will be set forth in part in the description which follows and in part will become apparent to those skilled in the art upon examination of the following or may be learned by practice of the invention. The objects and advantages of the invention may be realized and attained by means of the instrumentalities and combinations particularly pointed out in the appended claims."}
{"text": "Channel fading is a ubiquitous and fundamental characteristic of wireless communication systems. Fading deteriorates the link reliability of the wireless channel, thereby reducing system capacity and/or degrading user service experience. Diversity is a well-known principle that effectively combats wireless channel fading. The methods of achieving diversity include utilizing space, angle, polarization, frequency, time, and multipath. Diversity can be achieved at the transmitter and/or the receiver.\nConsider the forms of diversity that can be realized by using multiple antennas at the transmitter and/or receiver of a wireless communication system. These are grouped under the categories of transmit and receive diversity respectively.\nFIG. 1 illustrates a simplified diagram of a receiver 100 in a prior art wireless system that is equipped with multiple antennas (receive antenna 1 102, receive antenna N 102′) and exploits receive diversity. In this receiver 100, the multiple antennas (102. 102′) receive multiple versions of the same information-bearing signal. Assume that the wireless channel associated with any receive antenna is substantially statistically independent of the channel experienced by the other antennas. Then, the probability of each of the receive antennas being simultaneously faded is significantly smaller than that of any receive antenna being faded. Hence, the combined signal is much less likely to be faded, thereby improving the link reliability. In practice, the receive diversity gain is realized as follows. The signals received from the multiple receive antennas (102, 102′) are first individually processed with separate receive chains, each of which typically includes an analog signaling processing block (104, 104′), an analog-to-digital conversion block (106, 106′), and a digital signal processing block (108, 108′), respectively. The processed signals (109, 109′) are then combined in a combiner block 110. Combiner bock 110 may, for example, use selective combining or maximum ratio combining methods. The combiner 110 outputs signal 112, which may be subjected to further signal processing.\nSimilarly, FIG. 2 illustrates a simplified diagram of a transmitter 200 in a prior art wireless system equipped with multiple antennas (transmit antenna 1 202, transmit antenna N 202′) that exploits transmit diversity. In the FIG. 2 transmitter 200 illustrated here, the same information-bearing signal, source signal 204, is first split and pre-processed by splitter pre-processor 206 to generate multiple transmit signals (208, 208′), which are correlated with each other. These multiple transmit signals (208, 208′) are then individually passed through separate transmit chains including digital signal processing blocks (210, 210′), digital-to-analog conversion blocks (212, 212′), analog signal processing blocks (214, 214′) and transmitted with multiple antennas (202, 202′), respectively.\nTransmit diversity refers to the realization of diversity gain by sending multiple, correlated signals over a channel from the transmitter. Typically, transmit diversity techniques make use of multiple transmit antennas to transmit these correlated signals. Firstly, transmit diversity is not straightforward to realize, in general. Transmitting the same signal through multiple transmit antennas typically results in no diversity gain whatsoever.\nOne of the earliest transmit diversity techniques that was proposed is delay diversity, in which the transmitter sends multiple copies of the same information with different delays through different antennas. A more sophisticated version of this scheme which uses two transmit antennas was proposed by Alamouti described in S. M. Alamouti, “A simple transmitter diversity scheme for wireless communications,” IEEE Journal on Selected Areas in Communication, vol. 16, pp. 1451-1458, October 1998.\nLet the signal that is to be communicated be denoted by S(t) where t is assumed to be a discrete time instant. In the Alamouti scheme, two consecutive symbols are blocked off and transmitted over two time instants using the two antennas. Let X1(t) and X2(t) represent the output signals from the two antennas respectively, which may be expressed as\n      [                                                      X              1                        ⁡                          (              t              )                                                                          X              1                        ⁡                          (                              t                +                1                            )                                                                                      X              2                        ⁡                          (              t              )                                                                          X              2                        ⁡                          (                              t                +                1                            )                                            ]    =      [                                        S            ⁡                          (              t              )                                                            -                                          S                *                            ⁡                              (                                  t                  +                  1                                )                                                                                      S            ⁡                          (                              t                +                1                            )                                                                          S              *                        ⁡                          (              t              )                                            ]  \nSuppose that the time-varying channel responses from the two transmit antennas, e.g., two base station transmit antennas, to the receiver, e.g., a mobile receiver, are denoted by h1(t) and h2(t) respectively. For simplicity of explanation we can assume a flat channel but the more general case where the channel is frequency dependent can also be handled. If the channel coefficients are assumed to remain constant over two symbols, which is a mild assumption, the composite signal received by the mobile receiver can be represented byY(t)=h1X1(t)+h2X2(t)+W(t)Y(t+1)=h1X1(t+1)+h2X2(t+1)+W(t+1)which may be rewritten in terms of the original signal S(t) as\n            [                                                  Y              ⁡                              (                t                )                                                                                        Y              ⁡                              (                                  t                  +                  1                                )                                                        ]        =          [                                                                                    h                  1                                ⁢                                  S                  ⁡                                      (                    t                    )                                                              +                                                h                  2                                ⁢                                  S                  ⁡                                      (                                          t                      +                      1                                        )                                                              +                              W                ⁡                                  (                  t                  )                                                                                                                                          -                                      h                    1                                                  ⁢                                                      S                    *                                    ⁡                                      (                                          t                      +                      1                                        )                                                              +                                                h                  2                                ⁢                                                      S                    *                                    ⁡                                      (                    t                    )                                                              +                              W                ⁡                                  (                                      t                    +                    1                                    )                                                                        ]                  or      ⁢                          ⁢      alternatively        ,                  ⁢                  [                                                            Y                ⁡                                  (                  t                  )                                                                                                                          Y                  *                                ⁡                                  (                                      t                    +                    1                                    )                                                                    ]            =                                    [                                                                                h                    1                                                                                        h                    2                                                                                                                    h                    2                    *                                                                                        -                                          h                      1                      *                                                                                            ]                    ⁡                      [                                                                                S                    ⁡                                          (                      t                      )                                                                                                                                        S                    ⁡                                          (                                              t                        +                        1                                            )                                                                                            ]                          +                  [                                                                      W                  ⁡                                      (                    t                    )                                                                                                                                            W                    *                                    ⁡                                      (                                          t                      +                      1                                        )                                                                                ]                    \nIf the channel responses from the two transmit antennas to the receiver are known, it is straightforward to invert the transmitter code construction and extract the transmitted signal by the following transformation:\n                              [                                                                                          S                    ^                                    ⁡                                      (                    t                    )                                                                                                                                            S                    ^                                    ⁡                                      (                                          t                      +                      1                                        )                                                                                ]                =                              [                                                                                h                    1                    *                                                                                        h                    2                                                                                                                    -                                          h                      2                                                                                                            h                    1                                                                        ]                    ⁡                      [                                                                                Y                    ⁡                                          (                      t                      )                                                                                                                                        Y                    ⁡                                          (                                              t                        +                        1                                            )                                                                                            ]                                                  =                                            (                                                                                                              h                      1                                                                            2                                +                                                                                                h                      2                                                                            2                                            )                        ⁡                          [                                                                                          S                      ⁡                                              (                        t                        )                                                                                                                                                        -                                              S                        ⁡                                                  (                                                      t                            +                            1                                                    )                                                                                                                                ]                                +          noise                    which results in second-order diversity over a fading channel. The Alamouti scheme is simple, but requires the receiver to track the gains from each of the two transmit antennas separately, which normally requires two sets of pilots to be used. This is especially challenging in the cellular uplink, e.g., where a mobile device transmits to a base station receiver. Furthermore, the requirement of known transmit diversity techniques to use multiple transmitter chains, each of which normally includes both digital and analog signal processing blocks can be cost prohibitive in many applications.\nIn view of the above discussion, there is a need for improved methods and apparatus of achieving transmit and/or receive diversity in wireless communications systems. Methods and apparatus that achieve diversity while reducing the amount of signaling dedicated to pilots over known methods would be beneficial. Methods and apparatus that achieve diversity without the need for multiple transmit chains would also be beneficial."}
{"text": "4-Methyl-2-hydroxy-benzaldehyde is a known chemical entity [see e.g.: Ber. 11,773(1878) and Bull. Soc. Chim. France 35, 129-134(1924)], sofar however, its organoleptic properties have remained unrecognized.\nThough, it represents a simple higher homolog of prior used salicylic aldehyde, the compound of the present invention possesses unexpectedly an original flavor character which is distinct from that shown by this latter compound. In fact, salicylic aldehyde possesses a spicy and irritating odor similar to that shown by compounds such as benzaldehyde, acetophenone or nitrobenzene, accompanied by a marked phenolic note which renders its utilization rather difficult. In contradistinction therefrom, 4-methyl-2-hydroxy-benzaldehyde develops a sweet gustative note of caramel and burnt type reminiscent of some of the aspects of coumarin."}
{"text": "Color filters are widely used in projector system or virtual reality system to provide RGB primary colors from a white light source. The conventional projector system generally uses separately assembled color filters. More particularly, the color filters are separately manufactured and then assembled to the projector system. The manufacture tolerance is not important issue for the color filters.\nTo provide more compact color filters, method involving microelectronic and microlithography process are proposed to manufacture color filter films on transparent substrate.\nFIG. 1 shows a prior art process for color filter film (U.S. Pat. No. 5,711,889):                S100: preparing a substrate 112, as shown in FIG. 1A;        S102: forming a releasing layer 114 by evaporation atop the substrate 112, as shown in FIG. 1B;        S104: applying a photoresistor layer 116 atop the releasing layer 114 by conventional spin technology and then backing the photoresistor layer 116, as shown in FIG. 1C;        S106: exposing a predetermined portion of the photoresistor layer 116 by photolithography to define a dent 116′ on the photoresistor layer 116, where the optical filter will be positioned, as shown in FIG. 1D;        S108: etching a portion of the releasing layer 114 below the dent 116′, as shown in FIG. 1E;        S110: depositing an optical filter material 124 into the dent 116′ by electron gun evaporation, as shown in FIG. 1F;        S112: removing the remaining photoresistor layer 116 by suitable photoresistor remover such as acetone to expose underlying releasing layer 114, as shown in FIG. 1G;        S114: etching underlying releasing layer 114 to form optical filter material 124 with predetermined pattern, as shown in FIG. 1G.        \nThe optical filter material 124 with predetermined pattern can form Archimedes pattern on the substrate 112 and used with a light pipe to provide light filtering effect for a projector system.\nHowever, in above-mentioned process, the step of depositing an optical filter material 124 in step S110 is a “cold process” because the photoresistor layer 116 cannot resist high-temperature process. In conventional deposition by electronic gun, the substrate is heated to a relatively high temperature to enhance the quality of deposited film. Due to poor thermo stability of the photoresistor layer 116, optical filter material 124 with high quality is hard to form and extinct ratio of the optical filter using the optical filter material 124 is also poor."}
{"text": "Currently rolls of photographic film and paper are packaged in several different ways to facilitate room light loading into cooperating apparatus which dispenses or otherwise uses the film or paper. Commonly assigned U.S. Pat. Nos. 4,148,395 (Apr. 10, 1979, by Syracuse et al., and entitled \"Roll Package\") and 5,515,970 (May 14, 1996, by Ritchie et al., and entitled \"Light-Tight Package\") each discloses a package which includes a pair of flexible opaque end disks or covers attached to the ends of the core of the roll and an attached leader attached to the leading end of the length of film or paper forming the roll. A peripheral portion of each disk is folded over and adhered to the edge of an underlying convolution of the leader.\nCommonly assigned U.S. Pat. No. 5,133,171 (Jul. 28, 1992, by Chase et al., and entitled \"Light-Tight Packaging Method For Photosensitive Web Roll\") discloses a package that includes a pair of similar end disks and a leader. After a first convolution of the leader is wrapped onto the roll, peripheral portions of the end disks are folded over the first convolution. A second convolution is then wrapped onto the roll to capture the folded-over peripheral portions between the first and second convolutions and to make the package light-tight.\nAn existing practice for installing end disks on light-tight packages of the sort described above include manually grasping an end disk from a stack and then fastening it to the ends of the core. Another practice includes using a vacuum suction device to pick out a single end disk and then transfer it to the core of the roll. Each of the above practices, however, has numerous shortcomings that the present invention solves. In particular, an operator will invariable damage thin, flimsy, end disks when handling them manually. Moreover, multiple end disks are dispensed although only one is intended.\nAlthough not particularly related to the problem encountered by the inventors, U.S. Pat. No. 4,199,076 (Apr. 22, 1980, by Brown and entitled, \"Cup Dispenser\") discloses a dispenser of cups that has a slidable plate for supporting a stack of cups and an end plate having a cavity through which the cups are dispensed one at a time. The patent does not provide means for storing and dispensing thin, substantially flat objects.\nTherefore a need persists in the art for a method of storing and dispensing thin flimsy objects that is simple to use, protects the fragile objects from damage and is cost effective to manufacture."}
{"text": "The present invention relates to a Stirling-cycle refrigerator, and more particularly to a Stirling-cycle refrigerator capable of lowering a refrigerating temperature.\nStirling-cycle refrigerators comprise an expansion chamber defined by a displacer and a compression chamber defined by a piston, the expansion and compression chambers being interconnected through a heater or a high-temperature heat exchanger, a regenerator, and a cooler or a low-temperature heat exchanger. The displacer and the piston are operated out of phase with each other by a crank mechanism or a swash plate cam mechanism to cause a working gas to go through a continuous cycle of isothermal compression, isometric change, isothermal expansion, and isometric change for refrigeration.\nThe regenerator in the Stirling-cycle refrigerator of the type described contains a material made of metal or the like for storing heat. In general, when the temperature of a metal drops, the specific heat of the metal declines as does the heat storing ability thereof. With conventional regenerators filled with a metallic material, the lowest temperature which the refrigerator can achieve is governed by the specific heat of the metallic material used, and has been in the range of from 8.degree. to 9.degree. K."}
{"text": "Lithium ion batteries are secondary batteries having a structure in which, during charging, lithium dissolves out as ions from a positive electrode and migrates to the negative electrode to be stored therein; and on the contrary, during discharging, lithium ions return from the negative electrode to the positive electrode. Since lithium ion batteries have features such as high energy density and a long life cycle, lithium ion batteries are widely used as power supplies for domestic appliances such as video cameras; portable electronic devices such as laptop computers and mobile telephones; and electric tools such as power tools. Recently, lithium ion batteries are also applied to large-sized batteries that are mounted in electric vehicles (EV), hybrid electric vehicles, and the like.\nA lithium ion battery of this kind is configured to include a positive electrode, a negative electrode, and an ion conducting layer interposed between these two electrodes, and as this ion conducting layer, a separator formed from a porous film of polyethylene, polypropylene or the like, which is filled with a non-aqueous liquid electrolyte, is generally used. However, since such an organic liquid electrolyte which uses a flammable organic solvent as the solvent is used as an electrolyte, improvements in view of structure and material for preventing volatilization or leakage are needed, and installation of a safety device for suppressing temperature increase at the time of a short circuit and improvements in view of structure and material for preventing a short circuit are also needed.\nIn contrast, an all-solid lithium ion battery formed by solidifying the whole battery using a solid electrolyte that uses lithium sulfide (Li2S) or the like as a starting material, does not use a flammable organic solvent. Therefore, simplification of safety devices can be attempted, and the battery can be made as a battery which is excellent in terms of production cost or productivity. Also, the battery has a feature that the solid electrolyte can be laminated in series in a cell, and thus voltage increase can be promoted. Furthermore, in a solid electrolyte of this kind, since nothing but Li ion moves, side reactions caused by movement of anions do not occur, and it is expected that this leads to enhancement of safety and durability.\nA solid electrolyte used in such a battery is required to have high ionic conductivity as far as possible and to be stable chemically and electrochemically. For example, lithium halide, lithium nitride, lithium oxoate, and derivatives of these compounds are known as candidate materials for the solid electrolyte.\nIn regard to solid electrolytes of this kind, for example, Patent Document 1 discloses a sulfide-based solid electrolyte obtainable by incorporating a high temperature lithium ion conductive compound formed from lithium phosphate (Li3PO4), into a lithium ion conductive sulfide glass represented by general formula: Li2S—X (provided that X represents at least one sulfide among SiS2, GeS2, and B2S3).\nFurthermore, Patent Document 2 discloses, as a material that is crystalline and exhibits a very high ionic conductivity such as an ionic conductivity at room temperature of 6.49×10−5 Scm−1, a sulfide-based solid electrolyte characterized by including a lithium ion conductive substance as a composite compound represented by general formula: Li2S—GeS2—X (provided that X represents at least one of Ga2S3 and ZnS).\nPatent Document 3 discloses a lithium ion conductive sulfide ceramic having high lithium ion conductivity and a high decomposition voltage, the sulfide ceramic containing Li2S and P2S5 as main components and having a composition in which Li2S=82.5 to 92.5 and P2S5=7.5 to 17.5, as expressed in mol %, and preferably a composition (compositional formula: Li7PS6) in which Li2S/P2S5=7 as a molar ratio.\nPatent Document 4 discloses a lithium ion conductive material having a silver germanium sulfide mineral type crystal structure represented by chemical formula: Li+(12-n-x)Bn+X2−(6-x)Y−x (wherein Bn+ represents at least one selected from P, As, Ge, Ga, Sb, Si, Sn, Al, In, Ti, V, Nb, and Ta; X2− represents at least one selected from S, Se, and Te; and Y− represents at least one selected from F, Cl, Br, I, CN, OCN, SCN, and N3; while 0≦x≦2).\nPatent Document 5 discloses, as a solid compound that can be produced into a single layer in addition to the high fluidity of lithium ions, a lithium silver germanium sulfide mineral represented by general formula (I): Li+(12-n-x)Bn+X2−6-xY−x, and in this formula, Bn+ is selected from the group consisting of P, As, Ge, Ga, Sb, Sn, Al, In, Ti, V, Nb and Ta; X2− is selected from the group consisting of S, Se and Te; and Y− is selected from the group consisting of Cl, Br, I, F, CN, OCN, SCN, and N3, while 0≦x≦2."}
{"text": "This invention relates generally to the display of information intended for higher resolution display on a lower resolution format receiver.\nThere is increasing interest in displaying information normally displayed on a computer monitor on a television receiver. For example, set-top computers display computer information on an associated television receiver. In this way, conventional computer functions can be achieved without the cost of a monitor and in addition, computer functions can be applied to television reception. For example, set-top computers may implement electronic programming guides which allow access to programming information on the Internet or other sources. These guides also allow the user to select programs for viewing on his or her television receiver through interactions with graphical user interfaces.\nConventionally, computer systems display information on computer monitors in a non-interlaced or progressively scanned format making use of individual pixels which are displayed essentially simultaneously. Television receivers have conventionally utilized an interlaced display. An interlaced display may be divided into two interleaved segments. The odd numbered segments are scanned first and then the even numbered segments are interleaved in between to create a complete picture. Each of the segments, which are formed by a scanning electron gun, may be described as a “field.” A conventional television picture or frame may be made up of two fields. Recent digital television standards may also use progressive or non-interlaced scanning approaches.\nInformation, which was intended to be displayed on a computer monitor, may be degraded when displayed on a television screen. Among other things, the resolution on most television receivers is substantially less than that of computer monitors. The interlacing of the picture on the television receiver may also cause information to be lost. In addition, the color resolution of monitors may be better so that adverse hue and color saturation may occur when displaying computer information on television receivers.\nIn view of the lower resolution of most television receivers, a decimator may be used to intelligently remove horizontal rows of data from a frame meant for computer display. However, while this approach preserves the layout of the frame, it results in loss of resolution.\nIn some cases, content developed for display on a higher resolution display may actually be displayed on a lower resolution display than was intended. For example, content developed in a 1024×764 format may be displayed on a 640×480 display (such as a high definition TV display). Because the display resolution is lower, some of the data may be difficult to read.\nThus there is a continuing need for better ways to display information designed for a higher resolution format on a lower resolution display."}
{"text": "1. Technical Field\nThe present invention is in the field of sports racquets, especially tennis racquets, and more specifically relates to a new racquet component system with interchangeable components, which address performance characteristics.\n2. Background\nIn the field of tennis racquet design, much of the evolution in the design of racquets has been motivated concomitantly by a search for a larger sweet spot in the racquet and improved matching between the racquet specifications and an individual player's needs.\nThe sweet spot is defined as the most responsive part of the tennis racquet face. The sweet spot of the racquet provides improved power and control of racquet and ball, and transfers less vibration and shock to the racquet. From the player's perspective, the sweet spot is the area of the racquet face which gives the player the most power and control and the most solid hit. It is also the place on the racquet where the player feels the least vibration or shock. From a technical perspective, the sweet spot is the area of the racquet that is defined as having at least a 40% return rate. That is, for example, if a ball travels into the racquet at 100 miles per hour, it will leave the strings traveling at least 40 miles per hour.\nA number of changes have been made in racquet design in the prior art in pursuit of a larger sweet spot and greater dovetailing between the racquet specifications and individual player profiles. For example, the materials used to construct the racquet have gone from wood to metal, aluminum, graphite and various other combinations of synthetic and composite materials. The shape of the head of the racquet frame has been changed and made larger to increase the likelihood of ball contact and increase the area of the sweet spot. The larger head size of the racquet also increases the power of the racquet.\nThe new materials used to make the racquets are much lighter and stronger, giving the player more power and control but also causing the racquet to transfer more vibration to the player because of the lightness of the racquet. One response to this increase in vibration transferred to the player from the racquet has been to make the racquet stiffer. Another response to the problem of transferred vibration is seen in U.S. Pat. No. 4,983,242, Reed, which shows a tennis racquet that is formed with a damping sleeve to reduce the vibrations produced by the lighter racquet frames.\nIn addition to changes in the materials and the racquet head size, racquet design changes have also incorporated greater stiffness into the racquet. The increased stiffness in the racquet results in greater energy being returned to the ball, and again decreased vibration perceived by the player.\nWhat is still needed is a racquet which provides a quick and convenient means for custom tailoring a match between an individual player's needs and the racquet, allowing the player to change the performance characteristics of the racquet along with the \"feel\" of the racquet, for example, to increase the player's perception of the size of the sweet spot.\nWhat is also needed is a racquet component system that allows quick interchangeability of string heads. String heads can be interchanged to allow rapid repair of damaged strings during a game or match or at any time that it is not possible or appropriate to repair the damaged strings. Interchangeable string heads also allow for the rapid and convenient substitution of different string types and varying string tensions into the racquet. String heads may also be interchanged to allow the introduction of a different weight of string head into the racquet system.\nWhat is also needed is a racquet component system that allows quick and improved means for changing any component of the racquet component system to change the performance characteristics of the racquet under a long term or permanent change in conditions or the player's ability. For example, interchanging the racquet head or racquet frame to alter the control or power characteristics of the racquet or the damping and shock absorbing features of the racquet in response to player skill or to player injury.\nWhat is also needed is a racquet component system that allows quick interchangeability of racquet components to implement changing the performance characteristics of the racquet in response to changing conditions. These conditions encompass changes in either playing conditions, player conditions, racquet conditions or any combination of the three. Changes in the player's conditions include, for example, improvements in the player's skill and changes in the player's ability due to temporary injuries the player may sustain. Changes in playing conditions that a player can encounter include changes in the environment, such as the weather, or different court surfaces as well as varying styles of the player or the player's opponents.\nIt is, therefore, an object of the present invention to supply a racquet component system with interchangeable components, allowing better ability to control the \"feel\" of the racquet and a quick and convenient means for custom tailoring a match between an individual player and the racquet performance characteristics.\nIt is a further object of this invention to provide a racquet that allows quick interchangeability of string heads. String heads can be interchanged to allow rapid repair of damaged strings during a game or match or at any time that it is not possible or appropriate to repair the damaged strings. Interchangeable string heads also allow for the rapid and convenient substitution of different string types and varying string tensions into the racquet. String heads may also be interchanged to allow the introduction of a different weight of string head into the racquet system.\nIt is a further object of this invention to provide a racquet that allows quick and improved means for changing any component or any combination of the components of the racquet system, to change the performance characteristics of the racquet under constant conditions. For example, interchanging the racquet head or racquet frame to alter the control or power characteristics of the racquet or the damping and shock absorbing features of the racquet.\nIt is a further object of this invention to provide a racquet that allows quick interchangeability of racquet components to implement changing the performance characteristics of the racquet in response to changing conditions. These conditions encompass changes in racquet conditions, playing conditions, player conditions, or any combination of the three. Changes in the player's conditions include, for example, improvements in the player's skill and changes in the player's ability due to injuries. Changes in playing conditions that a player can encounter include changes in the environment, such as the weather, or different court surfaces as well as varying styles of the player's opponents, changes in racquet conditions would include, for example, changing the weight of the string head to change the damping characteristics of the racquet system or to change the overall weight of the racquet to effect a change of performance. An example of changing the overall weight of the racquet to effect a change of performance would be making the racquet light for playing close to the net versus making the racquet heavier for playing the baseline."}
{"text": "Spergualin is an antibiotic found by Umezawa (one of the present inventors) et. al. (U.S. Pat. No. 4,416,899). Since its discovery, many Spergualin-related compounds have been synthesized by Umezawa et. al. (U.S. Pat. No. 4,529,549, U.S. Pat. No. 4,556,735, etc.).\nThe compounds of the present invention have an immuno-modulating action with low toxicity and are expected as a medicine.\nConventional immuno-modulating agent generally have a strong side effect. Therefore, a compound having an excellent immuno-modulating action with reduced side effect and toxicity has been desired."}
{"text": "Oil has been a must for lubrication and cooling of an internal combustion engine. The quality of the oil plays an important roll for the service-life and the operation of the engine. Besides the quality of the oil, particles, especially metallic particle (basically ferrous metal particles), which are generated during the operation of the engine are also a factor that affects the service life of the engine. These metal particles are destructive to the engine, for they circulate with the oil to the moving parts of the engine and cause wearing and abrasion on the moving parts which in turn generate more ferrous metal particles. To overcome such a problem, an oil filter is devised to filter and clean the oil. The oil filter has been well known and being improved for may years. In general, an oil filter comprises a canister housing inside which filtering means, such as filtering paper or other porous material which allows the oil to flow through but stops the ferrous metal particles and thus cleans the oil. An example is disclosed in Taiwan patent publication No. 170,160.\nOil filters of this structure have been a standard for many years. However, due to the requirement of the minimum flow rate of the oil for lubrication and cooling of the engine, this oil filter structure is incapable to filter out minute particles, for it must have large enough openings to establish the required oil flow rate.\nTo overcome such a problem, magnet means has been disposed inside the oil filter to attract and thus remove the minute ferrous metal particles out of the oil flow. Some examples of the magnetic oil filters are demonstrated in U.S. Pat. Nos. 4,450,075, 4,613,435, 4,826,592, 4,851,116, 5,078,871 and 5,228,990 and Taiwan patent publication Nos. 165,334, 181,638 (which is corresponding to U.S. Pat. No. 5,228,990) and 195,840. All these patents comprise magnet means in different forms disposed inside the oil filter canister for removing ferrous metal particle out of oil flow within the canister.\nAlthough the conventional magnetic oil filters mentioned above do work well in removing minute metal particles from the oil flow, they have at least one major problem, namely, these oil filters comprise no metal particle collection means to allow the particles to be collected and securely held. Without the particle collection means, the ferrous metal particles that are attracted by the magnet means will gather around the magnet means and held thereon only by the magnetic force. The gathering of the metal particles may form a blockage to the oil flow.\nFurthermore, since the metal particles are only held by the magnetic force and since the magnetic force to some extent may be overcome by the vibration and shake of the engine or the automobile on which the engine is mounted during their operations, the ferrous metal particles that are attracted around the magnet means may re-enter the oil flow. This reduces the effectiveness of the conventional magnetic oil filters.\nAn even more serious problem caused by the metal particles that are shaken to release from the magnet means and re-enter the oilflow is that these particles may be magnetized to some degree by the magnet means during their contact engagement with the magnet means and thus in their re-circulation within the engine, they may attract each other and form a large particle which causes a more serious damage to the engine if allowed to circulate with the oil flow.\nMoreover, in some of the conventional magnetic oil filters, such as U.S. Pat. No. 5,228,990, the magnet means is fixed inside the canister by means of adhesive. Due to the high temperature condition during the normal operation of the engine, the oil flow that is thus maintained in a substantial high temperature deteriorates the adhesive and thus may cause failure of the adhesive. The failure of the adhesive results in un-securing the magnet means and thus damaging the oil filter. In certain conditions, the un-secured magnet means may also block the oil outlet and this causes an even more dangerous problem to the operation of the engine.\nIn some other magnetic oil filters, the magnetic means and the related parts are held inside the filter canister by means of frictional force provided by elastic deformation of the parts. However, due to the fact that to avoid unnecessary magnetization of the parts, some of the parts are made of aluminum which, as is generally well known, is less flexible and thus may not be able to provide desired elastic deformation for the provision of frictional force to hold the parts in position. As a result, the vibration and shake of the engine and/or the automobile may cause the magnet means and the related parts to disengage and thus damage the function of the oil filter."}
{"text": "Field effect transistors (FETs) are semiconductor devices used in a wide variety of electronics applications. A FET has three terminals: a source, a drain and a gate. During operation of the FET, current flows between source and drain terminals through a channel region. The gate electrode, positioned between the source and the drain, enables the current through the FET to be controlled based on the strength of the signal applied to the gate. The signal and bias present at the gate, source and drain determines the electric field profile in the channel region between the source and the drain. The performance of the FET, e.g., factors such as current gain, carrier mobility, and transconductance (gm), are determined by the profile of the electric field in the channel region.\nIn conventional FETs, the strength of the electric field varies over the length of the channel, being typically weaker near the source and stronger near the drain (in depletion mode). A non-uniform field can lead to decreased performance of the FET, because electrons near the source are accelerated slowly due to the relatively weak field in this region. Electrons near the drain may acquire too much energy due to the relatively strong field in this region, possibly causing damage to a gate insulator. An excessively strong electric field in one region can cause mobility degradation, hot electrons and impact ionization, and can generate gate leakage. FIG. 1A illustrates an example of a conventional metal semiconductor field effect transistor (MESFET) having a source 104, drain 106 and gate 108 formed on a substrate 102. FIG. 1B shows a curve 110 that illustrates an example of the magnitude of the electric field E in the channel region while the MESFET is in depletion mode. In this example, the magnitude of the electric field is relatively weak near the source and relatively strong near the drain. A non-uniform electric field, such as that illustrated by curve 110, can lead to decreased carrier mobility, non-linearity and non-constant transconductance.\nVarious techniques have been used to mitigate the effect of the non-uniform electric field, such as using a lightly-doped drain, delta doping of the channel, or using one or more field plates behind the gate. However, these methods lack flexibility to tailor the field in response to a range of operational voltages on the gate of the FET. Furthermore, no known field effect transistor provides constant transconductance."}
{"text": "1. Field of the Invention\nThe present invention relates to indication display units for vehicles, and more particularly to an indication display unit for vehicles which displays information concerning driving operation at a position corresponding to an instrument panel on a dash board.\n2. Description of the Prior Art\nConventionally, various information concerning driving operation such as the vehicle speed, engine revolutions, residual fuel quantity, time and other various warnings have been indicated by meters, a clock or alarm lamps provided on an instrument panel. Recently, however, in order to ensure sufficient space in a dash board, there has been proposed that a semi-transparent reflecting mirror is located at a position of an instrument panel, and an indicating means for indicating revolution frequency of engine, vehicle speed, residual fuel quantity, etc. to be essentially indicated is mounted on a lower surface of a meter hood, while various warning lamps which are not so essential are located behind the reflecting mirror. With this arrangement, the displayed image from the indicator is reflected on the reflecting mirror, and the images from the warning lamps are transmissively observed through the reflecting mirror (Japanese Utility Model Laid-open Publication No. 60-152544).\nIn this conventional device, both the transmissive image from the warning lamps and the displayed image from the indicating means are visually observed by the driver. Simultaneously, a sufficient space may be ensured in the dash board by not locating the indicating means behind the instrument panel, namely, in the dash board.\nHowever, as the meter hood is yet present, and the reflecting mirror is located at a position deeply inside the dash board, a sufficient space to be defined in the dash board is limited. Further, the meter hood is intended not only for the location of the meter indicator but also for prevention of reflection of the image from the warning lamps located behind the reflecting mirror on a windshield. Accordingly, even if the meter indicator is located at any position other than the meter hood, the meter hood is yet necessary to prevent the reflection of the image from the warning lamps on the windshield."}
{"text": "The need to seal between co-rotating rotor disks in the turbine or other sections of a gas turbine engine is a continuing problem. The environment which such seals must withstand includes exposure of at least one portion of the interdisk seal to the turbine working fluid having temperatures up to 2500F, withstanding the induced centrifugal force caused by high speed disk rotation, and accommodating thermal transient conditions caused by engine throttle changes. It is common for such seals to carry a portion of the rotating seal structure which is located between a stator vane assembly located axially intermediate the turbine disks. This rotating seal is typically comprised of an annular abradable member secured to the radially innermost portion of the stator vanes, and a series of knife edges extending circumferentially about the interdisk seal and extending radially outward into close contact with the abradable annular member. As will be well known to those skilled in the art, radial movement of the interdisk seal can cause the circumferential knife edges to move into contact with the abradable ring, opening a leakage path between the ring and knife edges during normal operating conditions.\nOne technique in the prior art to provide such interdisk seals has been the use of full ring seal members wherein a monolithic seal structure is established having a radially inner ring spanning the axial gap between the rotor disks, an annular outer wall member sealingly secured at axially opposite ends to the adjacent rotor disk rims, and a radially extending web member secured between the supporting ring and the outer wall for supporting the axially central portion of the outer wall. A drawback with this prior art seal member has been the occurrence of a thermal expansion mismatch between the rotor disks and the seal supporting ring.\nDuring periods of rapid engine acceleration, the working fluid temperature increases rapidly causing the disk rims and seal assembly to also increase in temperature. The interdisk seal, being of significantly lower mass than the disk members, increases in temperature more rapidly and hence experiences more rapid thermal expansion. The differential thermal expansion induced by the uneven temperature rise can cause excessive hoop stresses in portions of the full annular interdisk seal member which may be as great or greater than the hoop stress resulting from the rotation induced centrifugal force. Such monolithic interdisk seals must be fabricated of high strength materials in order to withstand the hoop stresses induced by rotation and the thermal growth mismatch discussed hereinabove. Such strength requires a heavier seal structure further penalizing the overall engine by imposing additional weight adjacent the rotor disk rims.\nWhat is needed is a lightweight interdisk seal which is able to accommodate differential thermal growth between the seal member and the adjacent turbine disks, which can withstand the rotation induced forces, and which is dimensionally stable in the radial direction to avoid excessive radial movement of the rotating knife edges."}
{"text": "1. Field of the Invention\nThe present invention relates to a transfer device for an automatic transfer switch, and particularly, to a transfer device for an automatic transfer switch, by which various electric components are allowed to be continuously supplied with stable power in a manner of blocking one air circuit breaker from being mechanically electrically closed when another air circuit breaker is closed due to a switching of states between the two air circuit breakers.\n2. Background of the Invention\nIn general, hospitals or telephone offices, in which power failure should not happen, use an automatic transfer switch (hereinafter, referred to as ATS) in order to ensure stable power supply. The ATS includes therein a regular air circuit breaker (hereinafter, referred to as a first ACB) and an emergency air circuit breaker (hereinafter, referred to as a second ACB). The first and second ACBs are mechanically connected to each other to be switched to each other. Accordingly, when a fault current is applied to one of the two ACBs, such fault current is detected. Consequently, a closed ACB is tripped and simultaneously another ACB is closed so as to be switched to each other. The switching of the ACBs allows a constant supply of stable power.\nUpon allowing the switching of the two ACBs for use, the two ACBs are typically connected to each other via a wire or a thin iron bar for allowing an interlocking of the ACBs.\nFIGS. 1 and 2 are perspective views showing an example of a related art device for interlocking two ACBs using a thin iron bar. Referring to FIGS. 1 and 2, the related art ATS includes frames 11 and 21 installed at the ACBs 1 and 2, respectively, closing links 12 and 22 rotatably coupled to the frames 11 and 21, respectively, to be rotatable by respective closing levers (not shown) installed within each ACB 1 and 2, trip links 13 and 23 rotatably installed in the frames 11 and 21, respectively, at one sides of the corresponding closing links 12 and 22 so as to cooperate with respective trip levers (not shown) installed in each ACB 1 and 2, and transmission rods 14 and 24 connected between the closing links 12 and 22 and the trip links 13 and 23 so as to transfer a rotational force. Here, the closing links 12 and 22 and the trip links 13 and 23 are installed independent of each other, and each of the transmission rods 14 and 24 is made of an integrally rigid material.\nWith the configuration of the related art ATS, when the first ACB 1 is closed, the first closing link 12 is rotated counterclockwise in the drawing due to the first closing lever installed in the first ACB 1. Responsive to this, the first transmission rod 14 is pulled upwardly such that the second trip link 23 of the second ACB 2 installed at a lower side in the drawing is moved upwardly. Here, an interlock lever is also moved upwardly by the second trip link 23, accordingly, the second trip lever installed in the second ACB 2 is pushed up such that the second ACB 2 is switched (transferred) to a tripped state.\nWhen the first closing link 12 is restored by a restoring spring (not shown), the pulled first transmission rod 14 is moved back to its original position and the second trip link 23 is also moved back to its original position, thereby releasing the tripped state of the second ACB 2.\nIn the meantime, a procedure that the second ACB 2 is switched to a closed state and the first ACB 1 to a tripped state is performed in an opposite manner to the aforesaid procedure.\nHowever, in the structure of the ATS of the related art, since the closing links 12 and 22 and the trip links 13 and 23 are independently installed, a problem is caused that a process of assembling the closing links 12 and 22 and the trip links 13 and 23 becomes complicated.\nAlso, a length of each transmission rod 14, 24 should be adjusted by cutting off the transmission rod 14, 24, which makes it difficult to fabricate the transmission rods 14 and 24 with exactly the same length between the closing links 12 and 22 and the trip links 13 and 23, and also makes it difficult to adjust an error occurred during the assembly operation, which may cause a mis-operation to occur due to a defective assembly.\nFurthermore, while each of the ACBs 1 and 2 is repeatedly closed and tripped, upon rotation of the closing links 12 and 22 and the trip links 13 and 23, an impact transferred to each transmission rod 14, 24 cannot be absorbed, thereby generating impact noise or degrading durability.\nIn addition, each transmission rod 14, 24 is formed of a thin rod, so it has a low intensity. Also, since each of the transmission rods 14, 24 has a structurally bent portion, it may have high possibility of being curved during operation, which may cause a mis-operation."}
{"text": "The present invention relates to a communication device, a communication system, and a communication method.\nRecently, measures are needed to prevent situations in which an attacker intrudes into the network of a critical infrastructure, such as a power plant, to take control of the systems of the infrastructure. The network of the critical infrastructure is protected by firewall devices and anti-virus software installed on terminals such as personal computers. In order to prevent the intrusion of attackers who intend to take over system control, stronger measures are needed. An example of a scheme that mitigates risks when an attacker intrudes into systems, there is a method that uses communication devices enabling whitelisting functions. The whitelisting function enabled by the communication device is a function that improves security level by registering information about authorized terminals included in authorized communications carried in networks on the communication device that can store whitelists and blocking unauthorized communications sent from terminals that are not authorized and not registered on the communication device that can store the whitelists.\nRegarding the background art of the present technique, there are Japanese Unexamined Patent Application Publication No. 2015-050767, Japanese Unexamined Patent Application Publication No. 2017-46149, and International Publication WO 2017073089.\nJapanese Unexamined Patent Application Publication No. 2015-050767 describes a whitelist-based network switch. The whitelist-based network switch includes a whitelist monitoring unit for storing a whitelist including authorized communication rules, monitoring one or more packets input through a plurality of switch interfaces based on the whitelist, and permitting communication of each packet conforming to the whitelist, and a whitelist management unit for updating the whitelist and transmitting an updated whitelist to the whitelist monitoring unit (See Abstract).\nJapanese Unexamined Patent Application Publication No. 2017-46149 describes a communication apparatus. The communication apparatus receives control information of first data and a plurality of types of header information of first data, the first data being received by a first data receiver; the apparatus selects a parameter from the plurality of types of header information of the first data based on a priority of a first data receiver group to which the first data receiver belongs and a storage condition, the priority being indicated by priority information, the storage condition indicating the number of entries of a whitelist that can be stored in a whitelist storage first memory; and the apparatus adds, to the whitelist, an entry that includes control information of the first data and at least one parameter selected above (See Abstract).\nInternational Publication WO 2017073089 describes a communication device, a system, and a method. A packet relay device is provided with packet reception units, a packet transfer unit, an S/W controller, packet transmission units, and an input-output interface, and a whitelist including authorized communication rules is automatically generated. For each data reception unit that receives data, it is possible to select controlling communication using the whitelist, or implementing data communication without using the white list (See Abstract).\nOn the other hand, Japanese Unexamined Patent Application Publication No. 2007-36374 introduces a security technique for networks using the Dynamic Host Configuration Protocol (hereinafter, DHCP). Japanese Unexamined Patent Application Publication No. 2007-36374 describes a packet transfer device, a communication network, and a packet transfer method. The packet transfer device has a plurality of ports, a protocol processing unit, and a controller. The packet transfer device transfers an Internet Protocol (IP) address allocation request sent from a client terminal in accordance with DHCP to a DHCP server, and receives an application for IP address allocation. In the transfer and reception, the packet transfer device stores information, which is an IP address and a Media Access Control (MAC) address, about the client terminal on the storage unit. The packet transfer device also stores information, which is an IP address and a MAC address extracted from an ARP packet, about the client terminal on the storage unit through an Address Resolution Protocol (ARP) resolution by the client terminal or through an ARP resolution by the packet transfer device. When the IP address of the stored DHCP packet is matched with the IP address of the stored ARP packet, the packet transfer device filters the port for the client terminal that has transmitted the ARP packet (See Abstract)."}
{"text": "This invention pertains generally to what is known in the wood products industry as the secondary breakdown of logs or lumber, and more specifically, to a lumber workpiece charger system that offers significantly improved feeding and handling of workpieces to an edger.\nIn the wood products field, market conditions in recent years have riveted attention on the need for increasingly precise processing control to maximize yield and to minimize waste. In addition, the same market conditions have focused concern on increasing the speed with which usable output can occur.\nOne area in which precision processing advances have certainly been seen in recent years is in the area of so-called edging of flitches, cants or boards preparatory to the making, for example, of dimensioned lumber. The edging process generally falls within what is referred to above as secondary breakdown activity.\nIn the past, there has been a somewhat unhappy mutual exclusivity between speed of handling in the edging process and precision in maximizing yield. For example, it is typical that, say, a board which is to be edged is transported laterally (normal to its longitudinal axis) beneath a scanning system which takes a look at the board's outline, taking into account wane where such exists, with scanning data then produced and used, as by a computer system, to determine how the edges of this board should be trimmed for maximum ultimate yield. A trimming decision might, for example, involve taking more wood off one long edge than off the other of the board, might require that the board be slightly angulated before feeding through an edger, as well as other things.\nThe so-called trimming decision regarding edging is realized in a variety of conventional pre-positioning or cuing devices, such as adjustable pins which rise into the travel path of a board to define what might be thought of as the proper \"leading edge\" orientation that should be used in shifting the board into the feed-intake station in an edger.\nMaintenance of proper position, from this point forward, requires careful positive handling of the board as the same as transported from its pre-positioned and angularly cued condition into the intake station of an edger, all to assure that when the board is \"taken\" by the edger, it will be guided along the edger's working axis with the proper predetermined angular and lateral orientation. However, equipment available in the past capable of accomplishing appropriate precision handling is relatively slow in operation. For example, prior art equipment is typically capable of handling successive work pieces at a rate approximating about 25- to about 30-pieces-per-minute under circumstances producing relatively gentle, non-damaging handling.\nEquipment offering high-speed transfer of a board from a pre-positioned cued condition to a handed-off condition to an edger typically has involved fairly violent jolting action which often offsets precision through failure to maintain proper board orientation.\nA general object of the present invention is to provide an improved edger charger system which offers the best of both worlds and the worst of neither. More specifically, an edger charger system is proposed herein which is capable of high-speed, anti-jolting, highly-precise handling of a board, or like workpiece, as the same is moved from a pre-positioned laterally and angularly cued condition toward ultimate hand-off for processing by an edger.\nAccording to a preferred embodiment of the invention, the proposed charger system includes a precharge station which is laterally offset from the intake station of an associated charger, which precharge station is adapted for receiving and holding, in succession, workpieces which are, based upon prescanning data, laterally and angularly cued for proper transport through the edger. The system further includes a positive-action pinch-roll mechanism which facially grips opposite faces of a workpiece held in the precharge station. This pinch-roll mechanism, moved under the influence of a computer-controlled, hydraulic, linear positioner, shuttles a gripped workpiece rapidly (i.e., with rapid acceleration, high speed, and rapid deceleration) to the intake station in the charger. This movement is accomplished without jolting, and as a consequence with little likelihood for disorientation of the properly pre-cued board.\nWhile the pinch-roll mechanism of the charger system continues to grip the board in the intake station, the usual clamp-roll mechanism in the edger closes upon the board, gripping the same, thus receiving a hand-off of the board before release of the latter by the charger's pinch-roll mechanism.\nAgain under the influence of the linear positioner, the pinch-roll mechanism is return-shuttled to the precharge station for similar handling of the next successive workpiece.\nPositive-action gripping of a workpiece from a pre-cued condition through gripped hand-off to an edger strongly promotes precision handling performance. Transfer motion produced under the influence of a computer-controlled hydraulic linear positioner enables rapid, anti-jolting handling, and thus adds significant handling speed without jeopardizing precision positioning.\nThese and other advantages which are offered by the system of the present invention will become more fully apparent as the description which now follows is read in conjunction with the accompanying drawings."}
{"text": "1. Technical Field\nThis invention generally relates to computer systems, and more specifically relates to apparatus and methods for determining the function of a computer system in response to user input.\n2. Background Art\nModern computer systems are tools that many people use on a daily basis. For many people, their time at work is spent to a large degree working on a computer system. Modern computer systems support the execution of many different software applications all at the same time. Sometimes when a user is actively using his or her computer system, some interrupting event occurs that interrupts what the user is doing. Pop-ups are a well-known example of an event that may interrupt the user's work. For example, if the user is typing in a work processor when a pop-up is activated, the pop-up typically is placed in front of the word processor screen. This requires a user to close down the pop-up window and then resume working in the word processor.\nPop-up blockers are known in the art, but are not 100% effective. In addition, there are may other types of events that could interrupt a user when actively working on their computer. For example, an e-mail program might provide a notification screen when a new e-mail is received. An antivirus program might provide a notification screen when a virus is detected. Other software applications may be running that interrupt the user periodically when user input is required. For example, a compiler could be executed to compile a sizable computer program, and the user could then perform other functions while the compiler does its work. However, the compiler, when finished, could interrupt the user. One can see from these simple examples that there exists a wide array of things that can interrupt a user while actively working on a computer. If a user is in the middle of entering text in a word processor, or entering data into a spreadsheet, or entering data into a database, or performing any other activity, the interruption is not only annoying, but negatively affects the user's efficiency and performance."}
{"text": "Copending U.S. application Ser. No. 09/723,778, filed concurrently herewith, entitled xe2x80x9cBallistic Aerosol Marking Process Employing Marking Material Comprising Vinyl Resin and Poly(3,4-ethylenedioxythiophene),xe2x80x9d with the named inventors Karen A. Moffat and Maria N. V. McDougall, now U.S. Pat. No. 6,383,561, the disclosure of which is totally incorporated herein by reference, discloses a process for depositing marking material onto a substrate which comprises (a) providing a propellant to a head structure, said head structure having at least one channel therein, said channel having an exit orifice with a width no larger than about 250 microns through which the propellant can flow, said propellant flowing through the channel to form thereby a propellant stream having kinetic energy, said channel directing the propellant stream toward the substrate, and (b) controllably introducing a particulate marking material into the propellant stream in the channel, wherein the kinetic energy of the propellant particle stream causes the particulate marking material to impact the substrate, and wherein the particulate marking material comprises toner particles which comprise a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), said toner particles having an average particle diameter of no more than about 10 microns and a particle size distribution of GSD equal to no more than about 1.25, wherein said toner particles are prepared by an emulsion aggregation process, said toner particles having an average bulk conductivity of at least about 10xe2x88x9211 Siemens per centimeter.\nCopending U.S. application Ser. No. 09/723,577, filed concurrently herewith, entitled xe2x80x9cBallistic Aerosol Marking Process Employing Marking Material Comprising Vinyl Resin and Poly(3,4-ethylenedioxypyrrole),xe2x80x9d with the named inventors Karen A. Motfat, Rina Carlini, Maria N. V. McDougall, and Paul J. Gerroir, now U.S. Pat. No. 6,467,871, the disclosure of which is totally incorporated herein by reference, discloses a process for depositing marking material onto a substrate which comprises (a) providing a propellant to a head structure, said head structure having at least one channel therein, said channel having an exit orifice with a width no larger than about 250 microns through which the propellant can flow, said propellant flowing through the channel to form thereby a propellant stream having kinetic energy, said channel directing the propellant stream toward the substrate, and (b) controllably introducing a particulate marking material into the propellant stream in the channel, wherein the kinetic energy of the propellant particle stream causes the particulate marking material to impact the substrate, and wherein the particulate marking material comprises toner particles which comprise a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxypyrrole), said toner particles having an average particle diameter of no more than about 10 microns and a particle size distribution of GSD equal to no more than about 1.25, wherein said toner particles are prepared by an emulsion aggregation process, said toner particles having an average bulk conductivily of at least about 10xe2x88x9211 Siemens per centimeter.\nCopending U.S. application Ser. No. 09/724,458, filed concurrently herewith, entitled xe2x80x9cloner Compositions Comprising Polythiophenes,xe2x80x9d with the named inventors Karen A. Moffat, Maria N. V. McDougall, Rina Carlini, Dan A. Hays, Jack T. LeStrange, and Paul J. Gerroir, now U.S. Pat. No. 6,503,678, the disclosure of which is totally incorporated herein by reference, discloses a toner comprising particles of a resin and an optional colorant, said toner particles having coated thereon a polythiophene. Another embodiment is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a resin and an optional colorant, said toner particles having coated thereon a polythiophene.\nCopending U.S. application Ser. No. 09/723,839, filed concurrently herewith, entitled xe2x80x9cToner Compositions Comprising Polypyrroles,xe2x80x9d with the named inventors Karen A. Moffat, Maria N. V. McDougall, Rina Carlini, Dan A. Hays, Jack T. LeStrange, and James R. Combes, now U.S. Pat. No. 6,492,082, the disclosure of which is totally incorporated herein by reference, discloses a toner comprising particles of a resin and an optional colorant, said toner particles having coated thereon a polypyrrole. Another embodiment is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a resin and an optional colorant, said toner particles having coated thereon a polypyrrole.\nCopending U.S. application Ser. No. 09/723,787, filed concurrently herewith, entitled xe2x80x9cBallistic Aerosol Marking Process Employing Marking Material Comprising Polyester Resin and Poly(3,4-ethylenedioxythiophene),xe2x80x9d with the named inventors Rina Carlini, Karen A. Moffat, Maria N. V. McDougall, and Danielle C. Boils, now U.S. Pat. No. 6,439,711, the disclosure of which is totally incorporated herein by reference, discloses a process for depositing marking material onto a substrate which comprises (a) providing a propellant to a head structure, said head structure having at least one channel therein, said channel having an exit orifice with a width no larger than about 250 microns through which the propellant can flow, said propellant flowing through the channel to form thereby a propellant stream having kinetic energy, said channel directing the propellant stream toward the substrate, and (b) controllably introducing a particulate marking material into the propellant stream in the channel, wherein the kinetic energy of the propellant particle stream causes the particulate marking material to impact the substrate, and wherein the particulate marking material comprises toner particles which comprise a polyester resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), said toner particles having an average particle diameter of no more than about 10 microns and a particle size distribution of GSD equal to no more than about 1.25, wherein said toner particles are prepared by an emulsion aggregation process, said toner particles having an average bulk conductivity of at least about 10xe2x88x9211 Siemens per centimeter.\nCopending U.S. application Ser. No. 09/723,834, filed concurrently herewith, entitled xe2x80x9cBallistic Aerosol Marking Process Employing Marking Material Comprising Polyester Resin and Poly(3,4-ethylenedioxypyrrole),xe2x80x9d with the named inventors Karen A. Moffat, Rina Carlini, and Maria N. V. McDougall, now U.S. Pat. No. 6,387,442, the disclosure of which is totally incorporated herein by reference, discloses a process for depositing marking material onto a substrate which comprises (a) providing a propellant to a head structure, said head structure having at least one channel therein, said channel having an exit orifice with a width no larger than about 250 microns through which the propellant can flow, said propellant flowing through the channel to form thereby a propellant stream having kinetic energy, said channel directing the propellant stream toward the substrate, and (b) controllably introducing a particulate marking material into the propellant stream in the channel, wherein the kinetic energy of the propellant particle stream causes the particulate marking material to impact the substrate, and wherein the particulate marking material comprises toner particles which comprise a polyester resin, an optional colorant, and poly(3,4-ethylenedioxypyrrole), said toner particles having an average particle diameter of no more than about 10 microns and a particle size distribution of GSD equal to no more than about 1.25, wherein said toner particles are prepared by an emulsion aggregation process, said toner particles having an average bulk conductivily of at least about 10xe2x88x9211 Siemens per centimeter.\nCopending U.S. application Ser. No. 09/724,064, filed concurrently herewith, entitled xe2x80x9cloner Compositions Comprising Polyester Resin and Poly(3,4-ethylenedioxythiophene),xe2x80x9d with the named inventors Karen A. Moffat, Rina Carlini, Maria N. V. McDougall, Dan A. Hays, and Jack T. LeStrange, the disclosure of which is totally incorporated herein by reference, discloses a toner comprising particles of a polyester resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), wherein said toner particles are prepared by an emulsion aggregation process. Another embodiment is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a polyester resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), wherein said toner particles are prepared by an emulsion aggregation process.\nCopending U.S. application Ser. No. 09/723,851, filed concurrently herewith, entitled xe2x80x9cloner Compositions Comprising Vinyl Resin and Poly(3,4-ethylenedioxypyrrole),xe2x80x9d with the named inventors Karen A. Motfat, Maria N. V. McDougall, Rina Carlini, Dan A. Hays, Jack T. LeStrange, and Paul J. Gerroir, now U.S. Pat. No. 6,485,874, the disclosure of which is totally incorporated herein by reference, discloses a toner comprising particles of a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxypyrrole), wherein said toner particles are prepared by an emulsion aggregation process. Another embodiment is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxypyrrole), wherein said toner particles are prepared by an emulsion aggregation process.\nCopending U.S. application Ser. No. 09/723,907, filed concurrently herewith, entitled xe2x80x9cToner Compositions Comprising Polyester Resin and Poly(3,4-ethylenedioxypyrrole),xe2x80x9d with the named inventors Karen A. Moffat, Rina Carlini, Maria N. V. McDougall, Dan A. Hays, and Jack T. LeStrange, now U.S. Pat. No. 6,387,581, the disclosure of which is totally incorporated herein by reference, discloses a toner comprising particles of a polyester resin, an optional colorant, and poly(3,4-ethylenedioxypyrrole), wherein said toner particles are prepared by an emulsion aggregation process. Another embodiment is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a polyester resin, an optional colorant, and poly(3,4-ethylenedioxypyrrole), wherein said toner particles are prepared by an emulsion aggregation process.\nCopending U.S. application Ser. No. 09/723,654, filed concurrently herewith, entitled xe2x80x9cProcess for Controlling Triboelectric Charging,xe2x80x9d with the named inventors Karen A. Moffat, Maria N. V. McDougall, and James R. Combes, now U.S. Pat. No. 6,365,318, the disclosure of which is totally incorporated herein by reference, discloses a process which comprises (a) dispersing into a solvent (i) toner particles comprising a resin and an optional colorant, and (ii) monomers selected from pyrroles, thiophenes, or mixtures thereof; and (b) causing, by exposure of the monomers to an oxidant, oxidative polymerization of the monomers onto the toner particles, wherein subsequent to polymerization, the toner particles are capable of being charged to a negative or positive polarity, and wherein the polarity is determined by the oxidant selected.\nCopending U.S. application Ser. No. 09/723,911, filed concurrently herewith, entitled xe2x80x9cloner Compositions Comprising Polyester Resin and Polypyrrole,xe2x80x9d with the named inventors James R. Combes, Karen A. Moffat, and Maria N. V. McDougall, the disclosure of which is totally incorporated herein by reference, discloses a toner comprising particles of a polyester resin, an optional colorant, and polypyrrole, wherein said toner particles are prepared by an emulsion aggregation process. Another embodiment is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a polyester resin, an optional colorant, and polypyrrole, wherein said toner particles are prepared by an emulsion aggregation process.\nCopending U.S. application Ser. No. 09/723,561, filed concurrently herewith, entitled xe2x80x9cElectrophotographic Development System With Induction Charged Toner,xe2x80x9d with the named inventors Dan A. Hays and Jack T. LeStrange, now U.S. Pat. No. 6,360,067, the disclosure of which is totally incorporated herein by reference, discloses an apparatus for developing a latent image recorded on an imaging surface, including a housing defining a reservoir storing a supply of developer material comprising conductive toner; a donor member for transporting toner on an outer surface of said donor member to a region in synchronous contact with the imaging surface; means for loading a toner layer onto a region of said outer surface of said donor member; means for induction charging said toner loaded on said donor member; means for conditioning toner layer; means for moving said donor member in synchronous contact with imaging member to detach toner from said region of said donor member for developing the latent image; and means for discharging and removing residual toner from said donor and returning said toner to the reservoir.\nCopending U.S. application Ser. No. 09/723,934, filed concurrently herewith, entitled xe2x80x9cElectrophotographic Development System With Induction Charged Toner,xe2x80x9d with the named inventors Dan A. Hays and Jack T. LeStrange, now U.S. Pat. No. 6,353,723, the disclosure of which is totally incorporated herein by reference, discloses a method of developing a latent image recorded or an image receiving member with marking particles, to form a developed image, including the steps of moving the surface of the image receiving member at a predetermined process speed; storing a supply of developer material comprising conductive toner in a reservoir; transporting developer material on a donor member to a development zone adjacent the image receiving member; and; inductive charging said toner layer onto said outer surface of said donor member prior to the development zone to a predefined charge level.\nCopending U.S. application Ser. No. 09/723.789, filed concurrently herewith, entitled xe2x80x9cElectrophotographic Development System With Custom Color Printing,xe2x80x9d with the named inventors Dan A. Hays and Jack T. LeStrange, now U.S. Pat. No. 6,463,239, the disclosure of which is totally incorporated herein by reference, discloses an apparatus for developing a latent image recorded on an imaging surface, including: a first developer unit for developing a portion of said latent image with a toner of custom color, said first developer including a housing defining a reservoir for storing a supply of developer material comprising conductive toner; a dispenser for dispensing toner of a first color and toner of a second color into said housing, said dispenser including means for mixing toner of said first color and toner of said second color together to form toner of said custom color; a donor member for transporting toner of said custom color on an outer surface of said donor member to a development zone; means for loading a toner layer of said custom color onto said outer surface of said donor member; and means for inductive charging said toner layer onto said outer surface of said donor member prior to the development zone to a predefine charge level; and a second developer unit for developing a remaining portion of said latent image with toner being substantial different than said toner of said custom color.\nThe present invention is directed to toners suitable for use in electrostatic imaging processes. More specifically, the present invention is directed to toner compositions that can be used in processes such as electrography, electrophotography, ionography, or the like, including processes wherein the toner particles are triboelectrically charged and processes wherein the toner particles are charged by a nonmagnetic inductive charging process. One embodiment of the present invention is directed to a toner comprising particles of a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), wherein said toner particles are prepared by an emulsion aggregation process. Another embodiment of the present invention is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), wherein said toner particles are prepared by an emulsion aggregation process.\nThe formation and development of images on the surface of photoconductive materials by electrostatic means is well known. The basic electrophotographic imaging process, as taught by C. F. Carlson in U.S. Pat. No. 2,297,691, entails placing a uniform electrostatic charge on a photoconductive insulating layer known as a photoconductor or photoreceptor, exposing the photoreceptor to a light and shadow image to dissipate the charge on the areas of the photoreceptor exposed to the light, and developing the resulting electrostatic latent image by depositing on the image a finely divided electroscopic material known as toner. Toner typically comprises a resin and a colorant. The toner will normally be attracted to those areas of the photoreceptor which retain a charge, thereby forming a toner image corresponding to the electrostatic latent image. This developed image may then be transferred to a substrate such as paper. The transferred image may subsequently be permanently affixed to the substrate by heat, pressure, a combination of heat and pressure, or other suitable fixing means such as solvent or overcoating treatment.\nAnother known process for forming electrostatic images is ionography. In ionographic imaging processes, a latent image is formed on a dielectric image receptor or electroreceptor by ion or electron deposition, as described, for example, in U.S. Pat. Nos. 3,564,556, 3,611,419, 4,240,084, 4,569,584, 2,919,171, 4,524,371, 4,619,515, 4,463,363, 4,254,424, 4,538,163, 4,409,604, 4,408,214, 4,365,549, 4,267,556, 4,160,257, and 4,155,093, the disclosures of each of which are totally incorporated herein by reference. Generally, the process entails application of charge in an image pattern with an ionographic or electron beam writing head to a dielectric receiver that retains the charged image. The image is subsequently developed with a developer capable of developing charge images.\nMany methods are known for applying the electroscopic particles to the electrostatic latent image to be developed. One development method, disclosed in U.S. Pat. No. 2,618,552, the disclosure of which is totally incorporated herein by reference, is known as cascade development. Another technique for developing electrostatic images is the magnetic brush process, disclosed in U.S. Pat. No. 2,874,063. This method entails the carrying of a developer material containing toner and magnetic carrier particles by a magnet. The magnetic field of the magnet causes alignment of the magnetic carriers in a brushlike configuration, and this xe2x80x9cmagnetic brushxe2x80x9d is brought into contact with the electrostatic image bearing surface of the photoreceptor. The toner particles are drawn from the brush to the electrostatic image by electrostatic attraction to the undischarged areas of the photoreceptor, and development of the image results. Other techniques, such as touchdown development, powder cloud development, and jumping development are known to be suitable for developing electrostatic latent images.\nPowder development systems normally fall into two classes: two component, in which the developer material comprises magnetic carrier granules having toner particles adhering triboelectrically thereto, and single component, which typically uses toner only. Toner particles are attracted to the latent image, forming a toner powder image. The operating latitude of a powder xerographic development system is determined to a great degree by the ease with which toner particles are supplied to an electrostatic image. Placing charge on the particles, to enable movement and imagewise development via electric fields, is most often accomplished with triboelectricity.\nThe electrostatic image in electrophotographic copying/printing systems is typically developed with a nonmagnetic, insulative toner that is charged by the phenomenon of triboelectricity. The triboelectric charging is obtained either by mixing the toner with larger carrier beads in a two component development system or by rubbing the toner between a blade and donor roll in a single component system.\nTriboelectricity is often not well understood and is often unpredictable because of a strong materials sensitivity. For example, the materials sensitivity causes difficulties in identifying a triboelectrically compatible set of color toners that can be blended for custom colors. Furthermore, to enable xe2x80x9coffsetxe2x80x9d print quality with powder-based electrophotographic development systems, small toner particles (about 5 micron diameter) are desired. Although the functionality of small, triboelectrically charged toner has been demonstrated, concerns remain regarding the long-term stability and reliability of such systems.\nIn addition, development systems which use triboelectricity to charge toner, whether they be two component (toner and carrier) or single component (toner only), tend to exhibit nonuniform distribution of charges on the surfaces of the toner particles. This nonuniform charge distribution results in high electrostatic adhesion because of localized high surface charge densities on the particles. Toner adhesion, especially in the development step, can limit performance by hindering toner release. As the toner particle size is reduced to enable higher image quality, the charge Q on a triboelectrically charged particle, and thus the removal force (F=QE) acting on the particle due to the development electric field E, will drop roughly in proportion to the particle surface area. On the other hand, the electrostatic adhesion forces for tribo-charged toner, which are dominated by charged regions on the particle at or near its points of contact with a surface, do not decrease as rapidly with decreasing size. This so-called xe2x80x9ccharge patchxe2x80x9d effect makes smaller, triboelectric charged particles much more difficult to develop and control.\nTo circumvent limitations associated with development systems based on triboelectrically charged toner, a non-tribo toner charging system can be desirable to enable a more stable development system with greater toner materials latitude. Conventional single component development (SCD) systems based on induction charging employ a magnetic loaded toner to suppress background deposition. If with such SCD systems one attempts to suppress background deposition by using an electric field of polarity opposite to that of the image electric field (as practiced with electrophotographic systems that use a triboelectric toner charging development system), toner of opposite polarity to the image toner will be induction charged and deposited in the background regions. To circumvent this problem, the electric field in the background regions is generally set to near zero. To prevent deposition of uncharged toner in the background regions, a magnetic material is included in the toner so that a magnetic force can be applied by the incorporation of magnets inside the development roll. This type of SCD system is frequently employed in printing apparatus that also include a transfuse process, since conductive (black) toner may not be efficiently transferred to paper with an electrostatic force if the relative humidity is high. Some printing apparatus that use an electron beam to form an electrostatic image on an electroreceptor also use a SCD system with conductive, magnetic (black) toner. For these apparatus, the toner is fixed to the paper with a cold high-pressure system. Unfortunately, the magnetic material in the toner for these printing systems precludes bright colors.\nPowder-based toning systems are desirable because they circumvent a need to manage and dispose of liquid vehicles used in several printing technologies including offset, thermal ink jet, liquid ink development, and the like. Although phase change inks do not have the liquid management and disposal issue, the preference that the ink have a sharp viscosity dependence on temperature can compromise the mechanical properties of the ink binder material when compared to heat/pressure fused powder toner images.\nTo achieve a document appearance comparable to that obtainable with offset printing, thin images are desired. Thin images can be achieved with a monolayer of small (about 5 micron) toner particles. With this toner particle size, images of desirable thinness can best be obtained with monolayer to sub-monolayer toner coverage. For low micro-noise images with sub-monolayer coverage, the toner preferably is in a nearly ordered array on a microscopic scale.\nTo date, no magnetic material has been formulated that does not have at least some unwanted light absorption. Consequently, a nonmagnetic toner is desirable to achieve the best color gamut in color imaging applications.\nFor a printing process using an induction toner charging mechanism, the toner should have a certain degree of conductivity. Induction charged conductive toner, however, can be difficult to transfer efficiently to paper by an electrostatic force if the relative humidity is high. Accordingly, it is generally preferred for the toner to be rheologically transferred to the (heated) paper.\nA marking process that enables high-speed printing also has considerable value.\nElectrically conductive toner particles are also useful in imaging processes such as those described in, for example, U.S. Pat. Nos. 3,639,245, 3,563,734, European Patent 0,441,426, French Patent 1,456,993, and United Kingdom Patent 1,406,983, the disclosures of each of which are totally incorporated herein by reference.\nU.S. Pat. No. 5,834,080 (Mort et al.), the disclosure of which is totally incorporated herein by reference, discloses controllably conductive polymer compositions that may be used in electrophotographic imaging developing systems, such as scavengeless or hybrid scavengeless systems or liquid image development systems. The conductive polymer compositions includes a charge-transporting material (particularly a charge-transporting, thiophene-containing polymer or an inert elastomeric polymer, such as a butadiene- or isoprene-based copolymer or an aromatic polyether-based polyurethane elastomer, that additionally comprises charge transport molecules) and a dopant capable of accepting electrons from the charge-transporting material. The invention also relates to an electrophotographic printing machine, a developing apparatus, and a coated transport member, an intermediate transfer belt, and a hybrid compliant photoreceptor comprising a composition of the invention.\nU.S. Pat. No. 5,853,906 (Hsieh), the disclosure of which is totally incorporated herein by reference, discloses a conductive coating comprising an oxidized oligomer salt, a charge transport component, and a polymer binder, for example, a conductive coating comprising an oxidized tetratolyidiamine salt of the formula \na charge transport component, and a polymer binder, wherein Xxe2x88x92 is a monovalent anion.\nU.S. Pat. No. 5,457,001 (Van Ritter), the disclosure of which is totally incorporated herein by reference, discloses an electrically conductive toner powder, the separate particles of which contain thermoplastic resin, additives conventional in toner powders, such as coloring constituents and possibly magnetically attractable material, and an electrically conductive protonized polyaniline complex, the protonized polyaniline complex preferably having an electrical conductivity of at least 1 S/cm, the conductive complex being distributed over the volume of the toner particles or present in a polymer-matrix at the surface of the toner particles.\nU.S. Pat. No. 5,202,211 (Vercoulen et al.), the disclosure of which is totally incorporated herein by reference, discloses a toner powder comprising toner particles which carry on their surface and/or in an edge zone close to the surface fine particles of electrically conductive material consisting of fluorine-doped tin oxide. The fluorine-doped tin oxide particles have a primary particle size of less than 0.2 micron and a specific electrical resistance of at most 50 ohms.meter. The fluorine content of the tin oxide is less than 10 percent by weight, and preferably is from 1 to 5 percent by weight.\nU.S. Pat. No. 5,035,926 (Jonas et al.), the disclosure of which is totally incorporated herein by reference, discloses new polythiophenes containing structural units of the formula \nin which A denotes an optionally substituted C1-C4 alkylene radical, their preparation by oxidative polymerization of the corresponding thiophenes, and the use of the polythiophenes for imparting antistatic properties on substrates which only conduct electrical current poorly or not at all, in particular on plastic mouldings, and as electrode material for rechargeable batteries.\nWhile known compositions and processes are suitable for their intended purposes, a need remains for improved marking processes. In addition, a need remains for improved electrostatic imaging processes. Further, a need remains for toners that can be charged inductively and used to develop electrostatic latent images. Additionally, a need remains for toners that can be used to develop electrostatic latent images without the need for triboelectric charging of the toner with a carrier. There is also a need for toners that are sufficiently conductive to be employed in an inductive charging process without being magnetic. In addition, there is a need for conductive, nonmagnetic toners that enable controlled, stable, and predictable inductive charging. Further, there is a need for conductive, nonmagnetic, inductively chargeable toners that are available in a wide variety of colors. Additionally, there is a need for conductive, nonmagnetic, inductively chargeable toners that enable uniform development of electrostatic images. A need also remains for conductive, nonmagnetic, inductively chargeable toners that enable development of high quality full color and custom or highlight color images. In addition, a need remains for conductive, nonmagnetic, inductively chargeable toners that enable generation of transparent, light-transmissive color images. Further, a need remains for conductive, nonmagnetic, inductively chargeable toners that have relatively small average particle diameters (such as 10 microns or less). Additionally, a need remains for conductive, nonmagnetic, inductively chargeable toners that have relatively uniform size and narrow particle size distribution values. There is also a heed for toners suitable for use in printing apparatus that employ electron beam imaging processes. In addition, there is a need for toners suitable for use in printing apparatus that employ single component development imaging processes. Further, there is a need for conductive, nonmagnetic, inductively chargeable toners with desirably low melting temperatures. Additionally, there is a need for conductive, nonmagnetic, inductively chargeable toners with tunable gloss properties, wherein the same monomers can be used to generate toners that have different melt and gloss characteristics by varying polymer characteristics such as molecular weight (Mw, Mn, MWD, or the like) or crosslinking. There is also a need for conductive, nonmagnetic, inductively chargeable toners that can be prepared by relatively simple and inexpensive methods. In addition, there is a need for conductive, nonmagnetic, inductively chargeable toners with desirable glass transition temperatures for enabling efficient transfer of the toner from an intermediate transfer or transfuse member to a print substrate. Further, there is a need for conductive, nonmagnetic, inductively chargeable toners with desirable glass transition temperatures for enabling efficient transfer of the toner from a heated intermediate transfer or transfuse member to a print substrate. Additionally, there is a need for conductive, nonmagnetic, inductively chargeable toners that exhibit good fusing performance. A need also remains for conductive, nonmagnetic, inductively chargeable toners that form images with low toner pile heights, even for full color superimposed images. In addition, a need remains for conductive, nonmagnetic, inductively chargeable toners wherein the toner comprises a resin particle encapsulated with a conductive polymer, wherein the conductive polymer is chemically bound to the particle surface. Further, a need remains for conductive, nonmagnetic, inductively chargeable toners that comprise particles having tunable morphology in that the particle shape can be selected to be spherical, highly irregular, or the like. Additionally, a need remains for insulative, triboelectrically chargeable toners that are available in a wide variety of colors. There is also a need for insulative, triboelectrically chargeable toners that enable uniform development of electrostatic images. In addition, there is a need for insulative, triboelectrically chargeable toners that enable development of high quality full color and custom or highlight color images. Further, there is a need for insulative, triboelectrically chargeable toners that enable generation of transparent, light-transmissive color images. Additionally, there is a need for insulative, triboelectrically chargeable toners that have relatively small average particle diameters (such as 10 microns or less). A need also remains for insulative, triboelectrically chargeable toners that have relatively uniform size and narrow particle size distribution values. In addition, a need remains for insulative, triboelectrically chargeable toners with desirably low melting temperatures. Further, a need remains for insulative, triboelectrically chargeable toners with tunable gloss properties, wherein the same monomers can be used to generate toners that have different melt and gloss characteristics by varying polymer characteristics such as molecular weight (Mw, Mn, MWD, or the like) or crosslinking. Additionally, a need remains for insulative, triboelectrically chargeable toners that can be prepared by relatively simple and inexpensive methods. There is also a need for insulative, triboelectrically chargeable toners with desirable glass transition temperatures for enabling efficient transfer of the toner from an intermediate transfer or transfuse member to a print substrate. In addition, there is a need for insulative, triboelectrically chargeable toners with desirable glass transition temperatures for enabling efficient transfer of the toner from a heated intermediate transfer or transfuse member to a print substrate. Further, there is a need for insulative, triboelectrically chargeable toners that exhibit good fusing performance. Additionally, there is a need for insulative, triboelectrically chargeable toners that form images with low toner pile heights, even for full color superimposed images. A need also remains for insulative, triboelectrically chargeable toners wherein the toner comprises a resin particle encapsulated with a polymer, wherein the polymer is chemically bound to the particle surface. In addition, a need remains for insulative, triboelectrically chargeable toners that comprise particles having tunable morphology in that the particle shape can be selected to be spherical, highly irregular, or the like. Further, a need remains for insulative, triboelectrically chargeable toners that can be made to charge either positively or negatively, as desired, without varying the resin or colorant comprising the toner particles. Additionally, a need remains for insulative, triboelectrically chargeable toners that can be made to charge either positively or negatively, as desired, without the need to use or vary surface additives. There is also a need for both conductive, inductively chargeable toners and insulative, triboelectrically chargeable toners that enable production of toners of different colors that can reach the same equilibrium levels of charge, and that enable modification of toner color without affecting the charge of the toner; the sets of different colored toners thus prepared enable generation of high quality and uniform color images in color imaging processes.\nThe present invention is directed to a toner comprising particles of a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), wherein said toner particles are prepared by an emulsion aggregation process. Another embodiment of the present invention is directed to a process which comprises (a) generating an electrostatic latent image on an imaging member, and (b) developing the latent image by contacting the imaging member with charged toner particles comprising a vinyl resin, an optional colorant, and poly(3,4-ethylenedioxythiophene), wherein said toner particles are prepared by an emulsion aggregation process."}
{"text": "The present invention generally pertains to paper products and is particularly directed to providing a paper product that is secure both from xerographic copying and from removal from secure premises.\nA xerographic copier includes a semiconductor layer that conducts electricity upon exposure to light but behaves as an insulator in the dark. In accordance with the xerographic copying process, the semiconductor layer senses the image to be copied when a mirror image of the image to be copied is reflected onto the semiconductor layer by a high energy light within a predetermined portion of the light spectrum. Such portion includes ultraviolet light. Light sources commonly used in xerographic copiers include quartz (tungsten) halogen lamps having an operational range between 400 and 900 nm and xenon lamps having an operational range between 380 and 1900 nm.\nIn the copying process, first, the semiconductor layer is electrostatically charged. Then, a mirror image of the image to be copied is projected onto the semiconductor layer by reflecting high energy light off of an original paper containing a printed image to be copied. In the areas of the semiconductor layer that sense the reflected light, the electric charge is dissipated. However, the residual charge is retained in the areas of the semiconductor layer that do not sense the reflected light, as a result of the high energy light being absorbed by the print on the original paper that defines the image to be copied. The semiconductor layer is then dusted with an oppositely charged toner powder which adheres to the residually charged areas to form the mirror image on the semiconductor layer. The image is transferred as a reproduced true image onto a copy paper that is brought into contact with the semiconductor layer and electrostatically charged from the rear to attract the toner powder onto the copy paper. The toner powder is then fused to the copy paper by heat to provide a permanent copy of the reproduced image on the copy paper.\nPaper products that are more or less secure from copying by xerographic photocopiers are known. One such paper product is distributed by the Fine Paper Company of Canada under the trademark \"NOCOPI\". This paper product is a standard paper characterized by a heavily dyed coating that is so dark that images printed on the paper can be seen with only great difficulty. When an attempt is made to copy the image by xerographic copying, the copy paper is turned totally dark. Another such paper product is made by the Xerox Research Center in Canada. This paper product contains a light sensitive matrix that is combined with the cellulose of the paper to cause the word void to appear on the copy paper when an attempt is made to copy an image from the paper product by xerographic copying. One drawback to this paper product is that after it has been imaged a few times with a halogen lamp (such as contained in many xerographic copiers) the paper product changes color so that the images on the paper product become unreadable to the naked eye.\nThere are systems for preventing the removal of a paper product from secure premises. One such system is described in U.S. Pat No. 3,665,449 to Elder et al. Such a system has been used to prevent unauthorized removal of documents from a secure area. A marker that produces a unique signal in response to an interrogation signal when transported through an interrogation zone of an electronic article surveillance (EAS) system is affixed to the document. The unique signal consists of harmonics of the interrogation signal that are uniquely characteristic of the marker material so that they may be distinguished from harmonics produced by other materials in response to the interrogation signal."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, and particularly to a MOS transistor based on high breakdown voltage specs and a manufacturing method thereof.\n2. Description of the Related Art\nA conventional high voltage MOS transistor has a structure wherein low-density diffused layers overlap with a gate electrode underneath the gate electrode to relax an electric field under the gate electrode to thereby suppress the occurrence of hot carriers. A method of manufacturing the conventional high voltage MOS transistor will be explained below with an N type MOS transistor as an example with reference to a process sectional views of FIGS. 3(a)–3(c).\nAn insulating film 302 such as an oxide film is formed on a P type semiconductor substrate 301 by known oxidation or a known CVD technique. Next, a resist pattern 303 is formed thereon by a known photolithography technique. Thereafter, an N type impurity such as phosphorus ions are implanted at a dose of 6.0 E+12 cm−2 by a known ion implantation technique. Next, the N type impurity is activated by a known diffusion technique to form N type low-density diffused layers 304 which serve as layers for relaxing source and drain electric fields in the MOS transistor (see FIG. 3(a)).\nNext, the resist pattern 303 is removed and a polysilicon film is deposited on the insulating film 302 by the known CVD technique. Afterwards, the polysilicon film is patterned by using a known photolithography and etching technique to form a gate electrode 305 (see FIG. 3(b)). Incidentally, at this time, the gate electrode 305 is formed in a structure in which the low-density diffused layers overlap with the gate electrode under the gate electrode, in such a manner that the gate electrode 305 covers parts of the N type low-density diffused layers 304 through the insulating film 302 by about 1.5 μm.\nNext, a resist pattern is formed by the known photolithography technique. Thereafter, an N type impurity such as As is implanted at a dose of 1.0 E+15 cm−2 by the known ion implantation technique to form N type high-density diffused layers 306 for drawing source and drain electrodes of the MOS transistor (see FIG. 3(c)). Incidentally, at this time, the N type high-density diffused layers 306 are formed away from the gate electrode 305.\nSubsequently, a high voltage MOS transistor having a structure in which low-density diffused layers overlap with a gate electrode underneath the gate electrode, is formed via contact formation and wiring formation. Incidentally, both the contact formation and wiring formation are done using a known technique and are not shown in the figure.\nIn order to improve breakdown voltage characteristics of the MOS transistor, Japanese Unexamined Patent Publication No. Hei 9(1997)-205205 cited as a patent document describes a method of forming high density layers for drawing or withdrawing source and drain electrodes at arbitrary and uniform positions, whereas Japanese Unexamined Patent Publication No. 2002-289845 cited as a patent document describes a method of forming low and high density layers in desired regions on a self-alignment basis.\nHowever, the above-described patent documents make no mention of the fact that the gate electrode is caused to overlap with the low density layers in a desired length. The method of manufacturing the high voltage MOS transistor having such a structure that the low-density diffused layers for field relaxation and the gate electrode both formed as described in the prior art have overlapped, was accompanied by the problem that there was a need to form the gate electrode after the formation of the low-density diffused layers, and when the photolithography technique was used, there was a need to determine the dimensions of the portions where the low-density diffused layers and the gate electrode overlapped in consideration of allowances for alignment between pattering for forming the low-density diffused layers and patterning for forming the gate electrode, thereby causing interference with device's miniaturization."}
{"text": "This invention generally relates to a mechanism and a process for coating portions of articles have dissimilar external configurations. More specifically, the invention relates to the application of a thermoplastic layer to coat tapped holes in specialty articles such as metal stampings.\nIt is well known to apply a thermoplastic resin powder such as nylon to threaded articles to form a \"patch\" which retards disengagement of the patched fastener with a mating fastener, as shown, for example, in U.S. Pat. Nos. 3,787,222 and 3,858,262. Mechanisms have also been developed for applying a protective coating to standard internally threaded fasteners at relatively high production rates, such as disclosed in U.S. Pat. No. 4,888,214, the disclosure of which is incorporated by reference herein. Similarly, other mechanisms have been developed for applying coatings to both standard and non-standard fasteners at lower production rates, such as disclosed in U.S. Pat. Nos. 5,141,771 and 5,362,327, also incorporated by reference herein.\nA common characteristic of the fasteners described in the patents listed above is that they possess only a single tapped hole (e.g., nuts), and their external dimensions are maintained within close tolerances.\nThere are many specialty articles, such as metal stampings, that contain multiple tapped holes and possess a relatively large variation in their external dimensions. While the tapped or threaded holes of such specialty articles can also advantageously utilize the protective coating or patch described above, prior art coating or patch applicators, including the mechanisms described in the above-referenced prior art patents, typically utilize external dimensions to position the threaded holes with the centerline of corresponding spray nozzles. Unfortunately, this approach is hampered with most stampings, for example, whose external surfaces often have burrs and rough edges which preclude precise positioning from these edges. In the past, to accomplish the coating of articles with varying external dimensions, the coating has been applied to manually positioned articles, at correspondingly low production rates.\nIt would, therefore, be desirable to provide an automated process, and an automated mechanism, for applying protective coatings or patches to threaded holes in articles having varying external configurations, while also providing a corresponding increase in production rates."}
{"text": "The human T1R family taste receptors include hT1R1, 2, and 3. The T1Rs belong to class-C G protein coupled receptors, and each Class-C GPCR consists of a large N-terminal extracellular domain and a C-terminal 7-transmembrane domain. It is generally known that hT1R1 and hT1R3 form a heteromeric receptor that modulates umami taste transduction and which recognizes umami tastants, while hT1R2 and hT1R3 form a heteromeric receptor that modulates sweet taste transduction and which recognizes sweet tastants. Therefore it is known that the umami taste receptor (hT1R1/hT1R3) and the sweet taste receptor (hT1R2/hT1R3) share a common subunit hT1R3."}
{"text": "The present invention relates to a key assembly for the keyboards of electric typewriters and similar machines and, more particularly, to a key assembly containing a plurality of vertically guided keyboard elements which, upon actuation, act upon domed contact beads of a contact mat arranged below the keyboard elements, with the contact mat, in turn, resting on a printed circuit board.\nA keyboard utilizing a key assembly of the kind referred to above is disclosed in German patent application (DE-AS) 23 50 176. When used with an electronic pocket calculator, the keyboard elements are inserted into the upper part of the housing while the contact mat and the printed circuit board are arranged in the lower part of the housing. By a spring plate, the keyboard elements are pushed into their normal positions into which they are lifted by the beads of the contact mat.\nThis arrangement requires high dimensional accuracy in assigning the reference surfaces between the keyboard elements and the contact mat. Under certain circumstances, greater tolerances are responsible already within the relatively small key pad of an electronic pocket calculator, for noticeable differences in the key travel.\nWith regard to keyboards of electric typewriters and similar machines which, not only from the surface, but also from the number of keys are considerably more voluminous, it appears that the principle of the known type of construction cannot be applied thereto. Owing to the guide lines on the ergonomic design of keyboards, certain requirements are placed with a view to the key travel and the operating force, demanding an exact, narrow-tolerated assignment of the keyboard element to the respective contact bead.\nIt is, therefore, the object of the invention to provide a key assembly for electric typewriters and similar machines which may use a contact mat and meets the above requirements."}
{"text": "This invention relates to boxes of chipboard or the like, supplied folded flat for filling and closing by means of side tabs and a top having a tongue that slides into the box and is provided with notches for locking the top in closed position. The invention is especially directed to such boxes in which the side tabs, as well as the top flap interlock to provide a more secure closure.\nThe object of the invention is to improve the side tab interlocking means to largely avoid the commonly experienced problem of tearing the portions of the side tabs intended to lock them together as the box is being closed so that a large proportion of the boxes are rendered useless and are discarded.\nMore specifically, the invention provides an interlocking arrangement which can be readily assembled without straining the box material and which securely locks the entire top closure structure together."}
{"text": "1. Field of the Invention\nThe present invention relates to a solid-state imaging device, which is an example of a semiconductor device for detection of physical amount distribution, and an imaging apparatus. Specifically, the present invention relates to a signal acquisition technique suitable for application to a solid-state imaging device or the like using a semiconductor device for detection of physical amount distribution in which a plurality of unit constituent elements with sensitivity to electromagnetic waves, such as light or radiant rays, input from the outside are arranged and which can read physical amount distribution, which is converted into electrical signals by the unit constituent elements, as electrical signals. In particular, the present invention relates to a structure of using a high-sensitivity pixel.\n2. Description of the Related Art\nA physical amount distribution detection semiconductor device in which a plurality of unit constituent elements (for example, elements) with sensitivity to electromagnetic waves, such as light or radiant rays, input from the outside are arranged in a line shape or in a matrix is used in various fields.\nFor example, in a field of visual equipment, a CCD (charge coupled device) type or MOS (metal oxide semiconductor) type or CMOS (complementary metal oxide semiconductor) type solid-state imaging device that detects a change of light (an example of an electromagnetic wave), which is an example of physical amount, is used. These solid-state imaging devices read physical amount distribution, which is converted into electrical signals by unit constituent elements (pixels in the solid-state imaging device), as electrical signals.\nFor example, the solid-state imaging device generates and accumulates signal charges by detecting electromagnetic waves, such as light or radiant rays, input from the outside in a photodiode, which is a photoelectric conversion element (light receiving element; photosensor) provided in an imaging portion (pixel portion) of a device portion, and reads the accumulated signal charges (photoelectrons) as image information.\nIn the case of a configuration for acquiring a color image, it is a mainstream to use an image sensor in which a color filter, which allows only a component having a specific wavelength to be transmitted therethrough, is arranged in each pixel and a required color component is restored by using a set of plurality of pixels.\nSpecifically, colors are distinguished by a color arrangement in which color reduction filters corresponding to, for example, red (R), green (G), and blue (B) colors, which are three primary colors, are used as a set of color filters and three primary color light components transmitted through the filters are separately detected by providing a semiconductor layer that detects light below the color filters. In addition, it may also be considered to use an arrangement in which white (Y), red (R), green (G), and blue (B) for luminance signal acquisition are combined. All the arrangements described above are called Bayer arrangements.\nIn a solid-state imaging element based on a single plate color method, each pixel has only information on a single color component as described above. Accordingly, for example, demosaic processing for restoring a required color component in each pixel by performing interpolation processing using color information on surrounding pixels is performed (for example, refer to JP-A-04-088784).\nOn the other hand, in a digital still camera or a movie camera, an improvement in quality of an image photographed under low illuminance has been an important issue in recent years. In the case of photographing an image under low illuminance, it is common to lower a shutter speed, to use a lens having a bright aperture value, or to use an external light source for visible light, such as a flash.\nIn this case, lowering the shutter speed causes hand trembling or blurring of a photographic subject. Moreover, in general, there is also a limitation in an aperture value of a lens. For this reason, it is difficult to realize the brightness more than a certain extent. In addition, in the case when the external light source for visible light is used, there is a problem that an atmosphere based on illumination in the place may be spoiled.\nSince a color temperature is low under a low illuminance condition, a light source that emits a large amount of infrared light is used in many cases. In addition, if invisible light, such as infrared light, is used as auxiliary light, the atmosphere is hardly spoiled. Therefore, a technique of increasing effective photographing sensitivity under a light source including a large amount of invisible light, such as infrared light, has been demanded.\nAs an example of measures against the problems described above, it is considered to add white (Y) for high-sensitivity luminance signals in a known color coding arrangement which is basically used. For example, JP-A-04-088784 discloses signal processing for obtaining a high-sensitivity image using an imaging element to which a YRGB arrangement, in which white (Y) for luminance signals is added in the Bayer arrangement where red (R), green (G), and blue (B) are arranged in a square matrix, is applied. In this structure, signal processing for obtaining high resolution and high sensitivity is explained by using a white pixel of pixels arranged in a checker pattern assuming that a color filter array, in which white pixels are arranged in the checker pattern, is applied and none of the pixels sense components (for example, infrared light) other than visible light components.\nThat is, in the arrangement proposed in JP-A-04-088784, a Y pixel in the checker arrangement has sensitivity to almost the entire visible light. Accordingly, a signal larger than that in a case of a configuration where green (G) pixels are arranged in a checker pattern is obtained. For this reason, it is possible to obtain a satisfactory S/N ratio of signals in the checker arrangement, which influences resolution or luminance information, compared with a case of a checker arrangement of green pixels."}
{"text": "This invention relates generally to gas turbine engines and more particularly to variable stator vane assemblies for use in such engines.\nGas turbine engines operate by combusting a fuel source in compressed air to create heated gases with increased pressure and density. The heated gases are ultimately forced through an exhaust nozzle, which is used to step up the velocity of the exiting gases and in-turn produce thrust for driving an aircraft. The heated air is also used to drive a turbine for rotating a fan to provide air to a compressor section of the gas turbine engine. Additionally, the heated gases are used for driving rotor blades inside the compressor section, which provides the compressed air used during combustion. The compressor section of a gas turbine engine typically comprises a series of rotor blade and stator vane stages. At each stage, rotating blades push air past the stationary vanes. Each rotor/stator stage increases the pressure and density of the air. Stators serve two purposes: they convert the kinetic energy of the air into pressure, and they redirect the trajectory of the air coming off the rotors for flow into the next compressor stage.\nThe speed range of an aircraft powered by a gas turbine engine is directly related to the level of air pressure generated in the compressor section. For different aircraft speeds, the velocity of the airflow through the gas turbine engine varies. Thus, the incidence of the air onto rotor blades of subsequent compressor stages differs at different aircraft speeds. One way of achieving more efficient performance of the gas turbine engine over the entire speed range, especially at high speed/high pressure ranges, is to use variable stator vanes which can optimize the incidence of the airflow onto subsequent compressor stage rotors.\nVariable stator vanes are typically circumferentially arranged between an outer diameter fan case and an inner diameter vane shroud. A synchronizing mechanism simultaneously rotates the individual stator vanes in response to an external actuation source.\nIn some situations, it is advantageous to divide the compressor section into upper and lower halves to expedite maintenance of the gas turbine engine. It is particularly advantageous, for example, in military applications when maintenance must be performed in remote locations where complete disassembly is imprudent. However, in dividing the compressor section into halves, the synchronizing mechanism must also be split apart. This creates two synchronizing mechanisms that must be actuated in unison to orchestrate simultaneous operation of all of the stator vanes. Synchronizing mechanisms that are located on the outer case can be accessed and spliced together easily. However, this is not the case for inner diameter synchronizing mechanisms, which cannot be accessed after assembly to attach the synchronizing mechanisms together. Thus, there is a need for an apparatus for coordinating actuation of split inner diameter synchronizing mechanisms."}
{"text": "The present disclosure relates to an apparatus for washing mechanical parts using a multipurpose aqueous- or solvent-based parts washer. Mechanical parts collect dirt, abrasion residue, used grease, and other debris during normal operation. During periodic maintenance extraordinary maintenance, repairs, and scheduled upgrades, mechanics disassemble parts from a larger mechanical element, such as a car engine or other industrial equipment. Individual parts and subassemblies must be washed before they are either discarded, diagnosed, and reinstalled or before they are reconditioned for further use.\nA parts washer is an apparatus that cleans parts, either individually or in groups of parts, including but not limited to machinery and machine parts. Parts washers can also clean elements such as chains, tools, and other elements susceptible to contamination from contact with greased or oiled parts. These cabinet-sized devices are an essential tool for any mechanic or worker who cleans parts in a workshop. For example, automobile mechanics place parts washers alongside tools or adjacent their work areas. The fundamental technology associated with parts washers is not unlike the technology associated with the cleaning of kitchen utensils and other food preparation accessories, the significant difference being that mechanical parts washer residue must be controlled before the effluents are released into the environment. Therefore, different cleaning solutions are often used, parts are generally washed infrequently once dirt is dried, oil-based effluents must be collected and confined, insoluble debris must be collected and filtered as sludge, and cleaning solutions are regenerated. The workshop environment in which parts washers are used also differs from location to location. Some parts washers use an aqueous cleaning solution to dissolve and remove grease, carbon, resin, tar, ink, and other debris. These parts washers use water, soap, and/or detergents, either commonly available or proprietary. Other more aggressive parts washers use hydrocarbon-based solvents or other solvents to degrease and wash parts. What is contemplated by this disclosure is a parts washer capable of using any type of cleaning solution, but more preferably, a parts washer capable of using either an aqueous-based or a solvent-based cleaning solution.\nBefore the arrival of parts washers, mechanics used small containers, collected a small volume of cleaning solution from a drum, and used ordinary sinks. Washing operations of the sink itself were required and a large volume of cleaning solution was wasted during each wash, since most cleaning solutions can be reused. A first generation of parts washers resemble a sink positioned over a reservoir where a cleaning solution is stored and recycled. An operator of the manual parts washer might push a pedal or take other action to activate a pump submerged in the cleaning solution reservoir and auxiliary heating element located within the reservoir might heat the cleaning solution to increase solubility of the circulating fluid, much like hot water is preferred during dishwashing operations. In a second generation of parts washers, manual washing operations were mostly been replaced with automated washing.\nThere are many advantages to manual parts washers over automated parts washing operations. For instance, they allow for tactile recognition of fine layers of dirt on parts having complex configurations. Manual cleaning also allows for the focus of cleaning efforts at a specific location, as well as to allow parts to be cleaned immediately. Other types of manual parts washers from the prior art include Safety-Kleen® Sink Models 16, 17, 30, 31, etc. Model 16, for example, comprises a metal sink with a pivoting lid placed on a drum that acts as the cleaning solution reservoir. These devices correspond to the first generation of parts washers.\nIn the 16/30, a pumping system is inserted within the drum and is functionally connected to a spray device in the sink. While highly effective over manual operations, this model can be functionally improved to enhance the overall effectiveness of parts washers. First, washing sinks are generally rectangular in shape for optimal usefulness, but these sinks are often nested on a cylindrical reservoir designed for rotational displacement ease and resistance to shock during transportation. As a result, the parts washer can become unstable if weight is placed in the corners of a rectangular sink placed atop a cylindrical reservoir. Cylindrical reservoirs, generally industrial drums, also have limited aesthetic value and often get dirty or bumped in work environments. If a reservoir is bumped, its top surface may be deformed, and a sink resting on the surface is then unstable.\nSecond, during the process of changing and replacing dirty cleaning solution that has been used to clean parts over a certain period of time, the sink must be removed from the reservoir and rested on a protected surface. The pump attachment connected to the bottom end of the sink is extracted from the drum and drips of cleaning solution. Parts washing also often requires the use of additional external light when they are placed in remote locations of work areas. Currently, an industrial lamp is attached to the lid at a location remote from the work area. The placement of a targeted spot lamp away from the washing area does not optimize illumination during the different washing operations.\nIn addition, since grease and oil exhibit under certain circumstances flammable properties that can be set aflame, fusible materials are used in the art to release holding devices of lids. The placement and orientation of fusible materials, when used in conjunction with locking mechanisms and latching mechanisms, is often troublesome. A fusible link must be in thermal proximity to a heat source but away from abrasion, shock, or work areas to optimize work performance of the overall parts washer.\nParts washers are generally stored where parts are removed or processed for convenient use. Auto repair shops, for example, are often equipped with multiple repair bays. A team of mechanics remove parts to be washed at different locations. Many mechanics prefer using only their assigned tools, though they are often less possessive of a parts washer. Immediate access to a proximate parts washer often trumps the need for an assigned washer. Confined spaces and other constraints associated with workshops warrant compact and portable devices. Parts washers must also be robust and durable under strenuous and prolonged use. Finally, industrial parts washers are designed with metal parts to be resistant to the shocks and impacts with the heavy metal parts cleaned within the sink. Metal, even when painted, can corrode if in contact with humid corrosive solutions, and once bumped, can change shape.\nWhat is needed is a third generation of parts washer capable exhibiting all the advantages of the first generation of washers, including but not limited to low cost, use of industrial drums of multiple sizes as a reservoir, lighting, portability, and simplicity of use with novel and useful features without losing the advantages of the first generation of parts washers. What is also needed is a series of operative and functional improvements to aid operators of the parts washer during washing."}
{"text": "Magnetic media are used in various electronic devices such as hard disk drives and magnetoresistive random access memory (MRAM) devices. Hard-disk drives are the storage medium of choice for computers and related devices. They are found in most desktop and laptop computers, and may also be found in a number of consumer electronic devices, such as media recorders and players, and instruments for collecting and recording data. Hard-disk drives are also deployed in arrays for network storage. MRAM devices are used in various non-volatile memory devices, such as flash drives and dynamic random access memory (DRAM) devices.\nMagnetic media devices store and retrieve information using magnetic fields. The disk in a hard-disk drive is configured with magnetic domains that are separately addressable by a magnetic head. The magnetic head moves into proximity with a magnetic domain and alters the magnetic properties of the domain to record information. To recover the recorded information, the magnetic head moves into proximity with the domain and detects the magnetic properties of the domain. The magnetic properties of the domain are generally interpreted as corresponding to one of two possible states, the “0” state and the “1” state. In this way, digital information may be recorded on the magnetic medium and recovered thereafter.\nMagnetic storage media typically comprise a non-magnetic glass, composite glass/ceramic, or metal substrate with a magnetically susceptible material between about 100 nm and about 1 μm thick deposited thereon by a deposition process, commonly a PVD or CVD process. In one process, a layer comprising cobalt and platinum is sputter deposited on a structural substrate to form a magnetically active layer. The magnetically susceptible layer is generally either deposited to form a pattern, or is patterned after deposition, such that the surface of the device has areas of magnetic susceptibility interspersed with areas of magnetic inactivity denominated by orientation of their quantum spin. Where domains with different spin orientations meet, there is a region referred to as a Bloch wall in which the spin orientation goes through a transition from the first orientation to the second. The width of this transition region limits the areal density of information storage because the Bloch wall occupies an increasing portion of the total magnetic domain.\nTo overcome the limit due to Bloch wall width in continuous magnetic thin films, the domains can be physically separated by a non-magnetic region (which can be narrower than the width of a Bloch wall in a continuous magnetic thin film). Conventional approaches to creating discrete magnetic and non-magnetic areas on a medium have focused on forming single bit magnetic domains that are completely separate from each other, either by depositing the magnetic domains as separate islands or by removing material from a continuous magnetic film to physically separate the magnetic domains. A patterned mask may be applied to a non-magnetic substrate, and a magnetic material deposited over exposed portions of the non-magnetic substrate, or the magnetic material may be deposited before masking and patterning, and then etched away in exposed portions. By one method, the non-magnetic substrate is topographically patterned by etching or scribing, and the magnetically susceptible material deposited by spin-coating or electroplating. The disk is then polished or planarized to expose the non-magnetic boundaries around the magnetic domains. In some cases, the magnetic material is deposited in a patterned way to form magnetic grains or dots separated by a non-magnetic area.\nSuch methods are expected to yield storage structures capable of supporting data density up to about 1 TB/in2, with individual domains having dimensions as small as 20 nm. All such methods typically result in significant surface roughness of the medium. Altering the topography of the substrate can become limiting because the read-write head of a typical hard-disk drive may fly as close as 2 nm from the surface of the disk. Thus, there is a need for a process or method of patterning magnetic media that has high resolution and does not alter the topography of the media, and an apparatus for performing the process or method efficiently for high volume manufacturing."}
{"text": "Quality of sleep is essential to good human health and well-being. Many types of beds have been developed, including mattresses and box springs, water beds, air mattresses, and memory foam mattresses, e.g. Mattresses and box springs comprising spring systems with coverings of pliable foam layers and possibly other materials have been in use for many years and are widely available. Many approaches and enhancements to beds and bedding have also been developed including mattress top coverings that include specially formed materials and foam to improve comfort. A wide range of materials, coverings, and configurations of blankets, sheets, pillows, and other bedding accessories are available.\nIn spite of a wide diversity of beds that are available, serious problems exist for substantially all types and configurations. Mattresses including spring systems sag over time and bugs (especially bed bugs), mites, human secretions, and other infestations and contaminants build up. Hence, many bed systems are not hygienic because there is no convenient way to launder, clean, or otherwise restore them. While a mattress with heavy coverings may be warm and cozy in the winter months, a lighter mattress covering that allows air to pass may be far more comfortable in the heat of summer. Conventional spring mattresses are also large and bulky, making them hard to store and move.\nWater beds are heavy, create damage if a leak or rupture occurs, and some persons have difficulty sleeping in them due to the motion of the water. Water beds normally include a heating system to warm the water in them that normally takes several hours to establish a requested temperature. Hence, while the water bed temperature is usually adjustable, a person may find a given temperature not comfortable and would need to wait several hours for the water bed to adjust. For at least the foregoing reasons, water beds are generally impractical for use in hotels.\nAir mattresses are normally not preferred for bedding as they don't allow air to pass, resulting in a warm and sweaty sleeping environment. In addition, some persons are unable to find comfortable sleeping positions on air mattresses.\nSome mattresses include viscoelastic foam. Viscoelasticity provides both viscous and elastic properties that enable a mattress to conform and support a person sleeping on it. The viscous nature of the material is sensitive to temperature and is normally configured so that the warmth of a person's body causes the viscoelastic foam to yield and conform; however, air does not easily pass through, making the mattress too warm for some persons, and viscoelastic mattresses are also susceptible to contaminations and infestations. Further, viscoelastic foam mattresses are generally heavy and also expensive relative to a conventional spring mattress.\nA key problem with all existing types of mattresses and beds is that while they may be tailored to persons of specific heights, weights, body shapes, and other personal aspects; it is impossible for them to be optimized for all persons. Hence, a mattress that provides stiffer springs in the region of a person's back for better back support of larger and heavier persons cannot help but generate an uncomfortably too stiff feeling for other persons who are smaller and lighter. Adjustable beds may provide some accommodations for certain sleeping positions, but if a person changes position the bed may need to be re-adjusted. Accordingly, adjustable beds substantially retain the issues with the various types of mattresses that they may be used with, and also tend to be heavy and expensive to produce."}
{"text": "This application relates to an apparatus for assembly of a press-fit work piece, such as a battery cell assembly.\nIn response to the demand for high-output battery applications such as for electric vehicles, modular battery cell stack assemblies have been developed. These stack assemblies may comprise multiple repeating units of a module that has generally planar battery cell members and other planar components such as gaskets, heat exchangers, and the like sandwiched between generally planar frame members. The planar frame members have retaining features thereon that are disposed perpendicular to the plane of the planar frame, which are adapted to engage retaining features on adjacent frame members in the stack when pressed together. The engagement of these retaining features holds the stack together and retains the sandwiched members securely between the frame members.\nOne way to assemble battery cell assemblies such as those described above has been to align the components and press them together in a manual operation, which required repeated manual applications of force to secure one frame to the next. In addition to being highly inefficient, such manual assembly can cause operator fatigue.\nAccordingly, the inventors herein have recognized a need for an apparatus for assembling a modular work piece such as a modular battery cell stack assembly that mitigates the above deficiency."}
{"text": "A variable reluctance (VR) motor is a brushless electric motor having a rotor and multiple pairs of inductive stator windings. Torque, in the direction of a point of minimum reluctance, is exerted upon the rotor in an amount proportional to the amount of current flowing through a pair of stator windings. Rotor torque, speed, and rotational direction are controlled in a VR motor by controlling the amount of current flowing through selected pairs of stator windings.\nVR motor control circuits which are constructed according to the prior art control the current flow through selected pairs of stator windings by alternating between a charge state and a discharge state. In the charge state, a voltage is impressed across the selected pair of windings until a desired current level is reached or slightly exceeded. A reverse voltage is then applied across the pair of stator windings to discharge the current flowing through the windings. In order to maintain a desired current through the windings, such a prior art control circuit must alternate between the charge and discharge states. The result is that the actual current is the average of the triangular current ripple which is created. This winding current ripple is highly disadvantageous because variations in rotor torque result and electromagnetic and acoustic noise is emitted.\nIn the illustrated preferred embodiment of the present invention a three state inductive load driver is used to control the current flow through each pair of stator windings in a VR motor. In a charge state the driver impresses a positive voltage across a pair of windings until a desired current level is reached. In a discharge state a reverse voltage is impressed across the windings until the current level has discharged to a desired level. An idle state, in which the windings are shorted, exists between the charge and discharge states. Since the current through a shorted ideal inductance is constant, the winding current during the idle state is constant except for decay due to minor resistive losses. Thus, a desired current level can be achieved and maintained with only a minimum of current ripple. Further, circuit complexity is minimized and the circuit operates at a low voltage and a high efficiency level."}
{"text": "One challenge posed by ink delivery systems for inkjet printers is air accumulation in the ink. When ink bubbles accumulate in the ink delivery system or in the print head, these bubbles can clog ink passageways and nozzles, thus harming print quality or preventing ink ejection altogether in at least part of the print head.\nAir accumulation via permeation is one mode by which air can accumulate in an inkjet ink delivery system. The print head ink-containing structure of an inkjet printer is typically a container made of lightweight polymer materials, which can be relatively permeable to air. Even where degassed ink is initially provided in the ink system, air can permeate through the polymer material of the ink reservoir wall over time, and dissolve into the ink. This dissolved air can produce bubbles and ultimately lead to failure of the print head."}
{"text": "Disposable gloves made from various synthetic and natural raw materials are used to protect against transmission of viral and bacterial agents and other pathogens, and to protect against chemical and radiological contamination. Users who must wear gloves for extended periods often suffer from dried or chapped hands, skin irritation, hand fatigue and premature wrinkling.\nGloves using aloe vera as a moisturizer are described in U.S. Pat. Nos. 6,274,154, 6,423,328, 6,630,152. Aloe vera is used in many skin care products, but it has an oil content that can have an adverse reaction when used with natural rubber and other materials, breaking down capability of rubber to act as a barrier. Some gloves also use lotions, which may have a similar effect on the barrier properties of gloves made from different raw materials. Other coatings used in gloves include cucumber extract and/or glycerin, but there is no evidence that common moisturizers used singly or in combination with another, have significantly improved the therapeutic properties of a disposable glove.\nWater soluble silicone has been used as a coating in disposable gloves since 1991, as a moisturizer and lubricant for ease of donning. Silicone has been shown to substantially reduce skin irritations when used in conjunction with natural rubber latex gloves.\nAllantoin, a botanical extract of comfrey also known as glyoxyldiureide has been recognized by the US FDA OTC panel as a skin protectant. It is believed to regenerate and stimulate cells, softens skin and acts to remove unhealthy tissue and is an anti-irritant.\nHydrolyzed collagen imparts protective colloid effect to formulations with anti-irritancy benefits. It increases the ability of skin to hold moisture. It is one of the long, fibrous structural proteins whose functions are quite different from those of globular proteins such as enzymes. Strong, tough bundles of collagen called collagen fibers are a major component of the extracellular matrix that supports most tissues and gives cells structure from the outside, but collagen is also found inside certain cells. Collagen has high tensile strength, and is the main component of cartilage, ligaments, tendons, bone and teeth. Along with soft keratin, it is responsible for skin strength and elasticity, and its degradation leads to wrinkles that accompany aging. It strengthens blood vessels and plays a role in tissue development."}
{"text": "Electronic distribution of financial product information from financial product originators to clients is a common service supplied by financial services providers. Typically, the service provider acts as a conduit for passing information related to various products between the financial content provider and clients. Such information includes new product offerings, updates to product attributes, such as a current price, and perhaps changes in the manner in which the product is classified.\nIt is not uncommon for a service provider to need to process dozens of data streams, each of which may be carrying information related to a large number of products. For example, a single data stream containing information about U.S. Treasuries typically contains data events related to several hundred different products.\nThere are various techniques for distributing raw financial services data to the clients. However, present techniques are labor or computationally expensive, treating the providers on an ad-hoc basis. In particular, in order to receive and process informational content so it can be distributed to downstream clients, the host needed to be hardwired to the specific structure and format of the incoming data stream. As a result, separate systems have developed for processing data from different sources. For example, one system may be dedicated to processing raw treasuries data and another independent system having its own structure may be dedicated to processing raw data related to corporate offerings. While acceptable for some operations, such ad-hoc techniques do not scale well and are difficult to integrate into a single comprehensive environment.\nIn addition, while present systems may be suitable for routing data from its source to various interested clients, conventional clients have been hard-wired to the host system with high-speed data connections. The proliferation of portable and wireless computer devices with low bandwidths and intermittent connectivity, along with the greater presence of smaller clients (both professional and individual) increases the need to provide a robust and adaptable data stream which is not susceptible to the inconsistencies inherent in communicating over a large communications network, such as the Internet.\nAccordingly, it is an object of the present invention to provide an integrated infrastructure product which is designed to process real-time content streams from a host system such that it can be distributed to authorized subscribers.\nIt is a further object to a system configured to process multiple streams of raw data in real time and providing a formatted and validated output which is suitable for down stream processing and distribution to clients.\nYet another object of the present invention is to provide a system which is flexible and easily adapted to permit the management of multiple streams of data from various sources and allow new streams and product types to be easily added or modified.\nYet a further object of the present invention is to provide validated and formatted offer data in a format which minimizes the amount of data required to transmit the offer to downstream services.\nYet another object of the present invention is to provide a data delivery mechanism which dynamically adjusts transmission rates to the speed at which a client can receive and process data while ensuring that updated data received by a client is current."}
{"text": "1. Field of the Invention\nThe present invention relates to a polymer scale preventive agent useful in polymerization of a monomer having an ethylenically unsaturated double bond, and a process of producing a polymer using said agent.\n2. Description of the Prior Art\nAs a method of polymerizing a monomer having an ethylenically unsaturated double bond, there have hitherto been known suspension polymerization, emulsion polymerization, solution polymerization, gas phase polymerization and bulk polymerization. In any type of the polymerizations, polymer scale is liable to be deposited on the areas with which the monomer comes into contact during polymerization, such as inner walls, agitation equipment, etc. of a polymerization vessel.\nDeposition of the polymer scale results in disadvantages that the yield of the polymer and cooling capacity of the polymerization vessel lowers; that the polymer scale may peel off and mix into the resulting polymer, thereby impairing the quality of formed products obtained by forming the polymer; and that the removal work of such a polymer scale not only is extremely laborious and hence time-consuming, but also may cause physical disorders, which are serious problems among the operators in recent years, due to unracted monomers which are contained in the polymer scales.\nHeretofore, with regard to preventing the polymer scale deposition on inner walls, etc. of polymerization vessels during polymerization of a monomer having an ethylenically unsaturated double bond, there is known a method of coating a polymer scale preventive agent consisting of an organic polar compound, such as an amine compound, a quinone compound or an aldehyde compound, on the inner walls, etc. of the polymerization vessel to form a coating film, and a method of adding these compounds into an aqueous medium where suspension polymerization is effected, as locally practiced, for example, in suspension polymerization of vinyl chloride [Japanese Patent Publication (KOKOKU) No. 45-030343 (1970)].\nHowever, such methods have the disadvantages that although a polymer scale preventive effect is exhibited while polymerization is repeated for up to about 5 or 6 batches, the preventive effect decreases (is inferior in durability) when the numbers or times of repeated batches of polymerization exceed 5 or 6. The decrease in the preventive effect is remarkable especially when used a water-soluble catalyst in polymerization, so that these methods can not be satisfied from the industrial view.\nIn order to overcome the above disadvantages, there are proposed, for example, methods of coating a polymer scale preventive agent containing, as an effective component, such as a condensation product of an aromatic amine compound with an aromatic nitro compound [Japanese Patent Publication (KOKOKU) No. 60-30681], a reaction product of a phenolic compound with an aromatic aldehyde [Japanese Pre-examination Patent Publication (KOKAI) No. 57-192414], a reaction product of a polyhydric phenol with an aliphatic aldehyde [Japanese Patent Publication (KOKOKU) No. 57-502169], and a reaction product of 1-naphthole with formaldehyde [Japanese Patent Publication (KOKOKU) No. 01-31523], and the like.\nIn case where each of these polymer scale preventive agents is coated on the areas, with which the monomer comes into contact during polymerization, such as inner walls, etc. of a polymerization vessel to thereby form a coating film, polymer scale deposition on the areas in the liquid phase of the polymerization vessel does not occur, even if polymerization batches are repeated about 100 to 200 times. Similarly, also in case where said water-soluble catalyst is used, polymer scale deposition on the liquid phase areas is prevented.\nHowever, said polymer scale preventive agents have a drawback that polymer scale deposition is liable to occur at and in the vicinity of the interface between a gas phase and a liquid phase which is located at the upper of the polymerization vessel, even if a coating film of a polymer scale preventive agent containing, as the effective component, the condensation product of said aromatic amine compound with said aromatic nitro compound and the like is formed.\nOnce polymer scale is deposited at or in the vicinity of the interface between the gas phase and the liquid phase, the deposited scale will grow gradually as polymerization runs are repeated, and at last peel off from the inner wall, etc. to get mixed with the polymeric product. If the polymeric product containing the polymer scale is processed into formed products such as sheets or the like, the polymer scale causes increase in fish eyes in the formed products, thereby lowering seriously the quality of the products.\nMeanwhile, where polymeric products obtained by polymerization are processed into formed products such as sheets, etc., the formed products are required to have a high whiteness. When a polymeric product is formed into a sheet or the like without adding any coloring agent, the resulting formed product is more or less colored. Such coloration is called initial coloration, which is desired to be as low as possible. However, when a coating film containing, as the effective component, the condensation product of said aromatic amine compound with said aromatic nitro compound, etc. is formed, the coating film may be peeled or dissolved during polymerization and mix into polymeric products, the resulting formed products will exhibit a low whiteness, that is, an inferior initial whiteness."}
{"text": "Conventionally, a CMOS imager pixel includes a phototransistor or photodiode operating as a light-detecting element. In operation, e.g., the pixel photodiode is first reset to a reset voltage that places an electronic charge across the capacitance associated with the photodiode. Electronic charge produced by the photodiode when exposed to illumination then causes charge of the photodiode capacitance to dissipate in proportion to the incident illumination intensity. At the end of an exposure period, the change in photodiode capacitance charge is detected and the photodiode is reset again. The amount of light detected by the photodiode is computed as the difference between the reset voltage and the voltage corresponding to the final capacitance charge.\nReferring to FIG. 1, the operation of a prior art pixel 10 is described. FIG. 1 shows the schematic diagram of a standard active pixel. A photodiode 11 produces a current proportional to the incident light intensity. The resulting photo current is integrated on a charge-sensing capacitor 13. The charge-sensing capacitor 13 is typically reverse-biased PN junction capacitance associated with the photodiode 11 and other parasitic capacitance.\nA MOS transistor 15 operates as a source-follower that buffers the voltage on the capacitor 13 nondestructively to a column line 23. A row select MOS switch 17 activates the source-follower transistor 15 when the particular row is selected by connecting the column current source 25 to the source of the source-follower transistor 15.\nThere are two primary ways to reset an active pixel, using a “soft” reset or using a “hard” reset. When using a “soft” reset, the voltage at the gate 21 of the reset transistor 19 is raised to a voltage that is no higher than the threshold voltage of the reset transistor, VRTTH, above the drain voltage of the reset transistor, typically at VDD. Generally, the voltage at the gate 19 is raised to the same potential as its drain voltage, VDD.\nAs the capacitor 13 is charged by the current from the reset transistor 19, the voltage at the sense node 27 increases, decreasing the gate-to-source voltage of the reset transistor 19. This in turn decreases the current from the reset transistor 19, and the rate of voltage rise at the sense node 27 decreases. As the gate-to-source voltage of the reset transistor 19 approaches its threshold voltage VRTTH of reset transistor 19, the current through the reset transistor 19 becomes extremely low, and the voltage at the sense node 27 rises very slowly. The voltage at the sense node 27 approaches approximately (VDD−VRTTH) but it never reaches a steady state because the rate of the voltage change becomes ever so slower. Then, the voltage at the gate 21 is lowered typically to ground, completing the reset process. At this time, the sense node 27 is reset to approximately (VDD−VRTTH).\nIn hard reset, the gate voltage of the reset transistor 19 is raised to a voltage greater than the drain voltage of the reset transistor by at least VRTTH. Typically, the gate voltage of reset transistor 19 is raised to VDD while the drain voltage of the reset transistor is maintained at a reset voltage VRESET that is lower than (VDD−VRTTH). This drives the reset transistor 19 into the triode region, thereby causing it to behave like a resistor.\nThe reset transistor 19 and the sense capacitor 13 behave like an RC circuit, and the sense node voltage approaches VRESET with an RC time constant, τ=RONC, where RON is the ON resistance of the reset transistor 19 and C is the value of the sense capacitor 13. Since the sense capacitance is on the order of a few femtofarads and the ON resistance is a few tens of kohms, the time constant is on the order of only a nanosecond.\nThus, the sense node typically reaches the full steady-state value VRESET within a few nanoseconds, which is much shorter than typical reset period of many microseconds. Then, the voltage at the gate 21 is lowered typically to ground, completing the reset process. At this time, the sense node 27 is reset to approximately VRESET.\nIt is well known that by using a “soft” reset, one can realize a lower reset noise, by a factor of √{square root over (2)}, compared to when using a “hard” reset. Thus, it is desirable from signal-to-noise ratio and sensitivity point of view to use a “soft” reset. However, since the sense node never reaches a steady state value, the voltage of the sense node is actually reset to different voltages depending on the initial condition on the sense node. This leads to substantial image lag. Even with popular double sampling method, a significant amount of image lag remains, which gives a blurry picture of moving objects.\nFurthermore, it is well known that by using a “hard” reset, one can substantially eliminate the image lag because the voltage to which the pixel is reset is always VRESET. However, the disadvantage of using a “hard” reset is that higher reset noise is realized.\nTherefore it is desirable to provide a imaging reset methodology and/or circuitry for an imager that provide a resetting capability, wherein reset noise is significantly reduced, image lag is substantially eliminated, a desirable signal-to-noise ratio is realized, and/or the imager realizes a desirable sensitivity."}
{"text": "A router is an information technology device included in a given communication network, such as a packet based network, whose function it is, in particular, to provide a junction between the given network and at least one other network, to direct data packets through the given network, and to transmit data packets from the given network to another network.\nAn attack against a router can bring about a malfunction of the router and various detrimental consequences, such as breakdown of the network, an interruption of service, a disruption to the activities of the users of the network, and substantial cost and negative advertising.\nSuch attacks often target routers because of their critical role in a network, and they can be of different types, such as a “man in the middle” attack to monitor data traffic within the network, a denial of service attack, or a subversion attack.\nIn subversion attacks, the attacker can be connected to the router with current user rights or possibly with administrator rights. Such rights authorize the attacker to modify the router's configuration and to control all processes executed in the router.\nThere are various routing protocols, such as the OSPF (“Open Shortest Path First”), the BGP (“Border Gateway Protocol”), or the RIP (Routing Information Protocol”) protocols to direct data traffic through networks. During a subversion attack, the attacker can modify the process that generates messages of the routing protocol in order to use fake addresses and to direct data packets over bad routing. For example, data packets can be directed towards equipment belonging to the attacker so that the latter can access such data packets in order potentially to be able to modify them. Such bad routing information can be propagated to all routers in the network.\nTo ensure secure data exchanges between two routers of a network, certain mechanisms authenticate all messages generated by the routing protocol used by the network. This authentication is carried out by means of secret keys that are deployed and memorized in all routers of the network and which can be changed by each router's administrator. Such authentication mechanisms have the following disadvantages:                when an attacker is in possession of a router's administrator rights, he/she can modify the secret keys required for the router's security functions,        an attacker can potentially change the implementation of the routing protocol to modify the content of messages sent or received, and        secret key management is very onerous, to the extent that secret keys are generally never changed.        \nThere is therefore a need for a packet based network operator to prevent or limit the consequences of an attack against a network router, in which an attacker could direct all or part of data traffic to a corrupted network address."}
{"text": "Field of the Invention\nThe disclosed and claimed concept relates to a stator assembly for a progressing cavity pump and, more specifically, to a stator assembly wherein the helical passage is a flexible helical passage.\nBackground Information\nProgressing cavity pumps are often referred to as “Moineau” pumps, in recognition of their inventor, Rene Moineau, who obtained U.S. Pat. No. 1,892,217. Progressing cavity pumps are used in various industries to pump materials such as, but not limited to, viscous fluids, semi-solids, fluids with solids in suspension, and solids. Exemplary materials transported by a progressing cavity pump include, but are not limited to, oil, sewage, fracking fluids or the like. Generally, a progressing cavity pump (also known as a helical gear pump) includes an elongated rotor having one or more externally threaded helical lobes, or “splines,” rotatably disposed in a stator assembly or stator body defining a helical passage. In one embodiment, the helical passage includes one more lobes than the helical rotor. The elongated helical passage includes a plurality of helical grooves that form a plurality of cavities with the stator. As the rotor turns within the stator, the cavities progress from a suction end of the pump to a discharge end. In other embodiments, there are an equal number of rotor splines and stator lobes, but the rotor splines are sized and shaped so as to define cavities within the stator lobes. In an exemplary embodiment, each lobe of the rotor is, in theory, constantly in general contact with the stator at any transverse cross section; this has the effect of creating a plurality of empty spaces between the stator and the rotor. It is noted that the clearance, or interference, at a location wherein a rotor spline is not fully seated in a stator lobe, may be variable, i.e., less than substantial engagement. That is, for example, in an embodiment wherein a stator passage has an arcuate end surface and a linear lateral surface, it is desirable to ensure the rotor seals against the arcuate end surface of the stator; this ensures the cavity, and therefore the fluid therein, moves forward. It is desirable, but less important, that the rotor seals against the linear lateral surface of the stator.\nAs the rotor rotates, the empty spaces advance from the suction end of the helical passage to the discharge end of the helical passage. Further, the empty spaces are isolated from each other by the points of contact between the rotor and the stator, which are often referred to as “seal lines.” As the rotor rotates within the stator, the empty spaces “move” or progress with a helical motion along the length of the helical passage. In operation of a progressing cavity pump, the empty spaces are filled with a material that is to be moved. Thus, as the empty spaces progress, the material is moved from one end of the stator to the other end of the stator as the rotor rotates relative to the stator. Due to the shape and geometry of the stator and the rotor, the rotor will move laterally or precess relative to the stator as the rotor rotates within the stator. In other words, the rotor moves eccentrically relative to the stator in addition to rotating within the stator.\nIn an exemplary embodiment, shown in FIG. 1, a progressing cavity pump 1, includes an elongated helical rotor 2, and a stator assembly 3 defining an elongated helical passage 4. In the exemplary embodiment shown, the rotor has a single lobe and, therefore, has a generally circular cross-sectional shape. The helical passage (shown in cross-section) has an obround shape. As used herein, an “obround” shape includes opposed generally arcuate surfaces and opposed generally parallel, generally linear surfaces; what may be colloquially identified as a “pill” shape. In operation, the rotor 2 reciprocates between the two ends of the helical passage.\nTo ensure that the rotor is “constantly in substantial contact with the stator at any transverse cross section” the stator helical passage is typically lined with a resilient material, such as but not limited to an elastomeric material. That is, in an exemplary embodiment, the stator assembly includes a rigid support assembly defining the helical passage and the liner is disposed thereon. As the rotor rotates and reciprocates between the two ends of the helical passage, in the exemplary embodiment shown in FIG. 1, the resilient material is compressed between the rotor and the support structure. Further, if the material being moved is a fluid with suspended solids, the solids may pass between the resilient material and the rotor.\nThis configuration has several disadvantages including the degradable nature of the resilient material liner. That is, the compression of the resilient material liner causes rapid wear and tear on the liner leading to the need for replacement. As used herein, “rapid” degradation is a relative term; a resilient material degrades more rapidly than a durable material. Further, solids passing between the resilient material and the rotor also damage the resilient material liner. Also, the resilient material liner may react with, or be degraded by, the material being moved. Another disadvantage is that rigid stator assemblies are difficult and/or expensive to construct. That is, such stator assemblies are typically created by hydroforming, rolling a metal tube, cold drawing a metal tube, hot extrusion of a metal tube, boring a metal tube using a method such as, but not limited to, electrical discharge machining, and electroforming with metal deposition.\nIn another embodiment, not shown, the stator assembly is made substantially of a resilient material. While the resilient material may have a rigid outer housing, the helical structure and support is formed by the resilient material. This embodiment also allows for substantial constant contact between the rotor and the stator assembly, and, allows for solids to pass between the rotor and stator. This embodiment is, however, also subject to rapid degradation. Further, as the stator helical passage is generally resilient, the progressing cavity pump of this embodiment is limited to lower pressures and lower transfer speeds. That is, at a higher pressure, the stator will distort allowing back-flow of the material over the rotor.\nIn another embodiment, not shown, the stator assembly is made of a rigid material with no liner. Typically, both the rotor and the stator are made from a durable material, i.e., a non-resilient material. While a durable material is less subject to wear-and tear, the friction between the two durable material elements will cause wear-and-tear to both the rotor and the stator. Further, with rigid materials forming both the rotor and the stator, particles cannot pass therebetween. That is, a solid trapped between the rigid rotor and stator will be crushed causing additional wear and tear to the components. Alternatively, with a larger or more durable particle, the rotor will flex, possibly bending the rotor permanently. As such, and as used herein, a progressing cavity pump wherein a durable rotor engages, or moves over, a durable stator is a “self-damaging” progressing cavity pump. One solution to the issue with particles in a self-damaging progressing cavity pump is to allow for a small gap between the rotor and the stator; that is, the rotor and stator are not “constantly in contact.” This configuration, however, allows for back-flow of the material between adjacent cavities. That is, this configuration is less efficient. Further, in this embodiment, the stator is typically made by one of the expensive methods noted above.\nFurther, as noted in U.S. Pat. No. 8,905,733 there is an advantage to having turbulent flow of a fluid adjacent the stator surface within a progressing cavity pump. In that patent, the turbulent flow is created or enhanced by grooves in, for example, the surface of the stator helical passage. These grooves, however, must be machined into the stator helical passage surface either during the formation of helical passage or sometime thereafter. As such, the grooves are expensive to incorporate into the stator.\nIt is understood that a progressing cavity pump includes a drive assembly with a drive shaft that causes the rotor to rotate within the stator thereby creating the pump action. That is, a rotary motion is converted to a fluid action, i.e., pumping. As is known, however, the rotor/stator assembly with minor geometric differences may have a fluid pumped therethrough thereby causing the rotor to rotate. That action is then transferred to the drive shaft and drive assembly. That is, a fluid motion is converted into a mechanical motion. Thus, it is understood that while the following discussion addresses a rotor/stator assembly as a pump, the same rotor/stator assembly may be used to create a rotational motion, i.e., may be used as a drive device, e.g., for a drill.\nThere is, therefore, the need for an improved progressing cavity pump wherein the components are not subject to rapid degradation, are not self-damaging, and do not allow for back flow of the material being transported."}
{"text": "The present invention relates to a system for displaying a multilayer image on a screen and, in particular, it relates to a technique for determining a request destination when a user points a position on the screen and makes a process request.\nHitherto, systems for displaying multilayer images are known. Such systems display a plurality of images through alphablending the images according to their respective transparency information and superposing them. The transparency information can be set for each point on the screen. Color information and the transparency information are combined to display one point on the screen. The transparency information set for each point on the screen is called an alpha value, which ranges from 0 (completely transparent) to 1 (opaque).\nProblems sometimes occur on a display screen that displays a multilayer image when a user selects an object located in one layer. For example, suppose a semitransparent enlarging or reducing button is displayed on a map. In this case, even if the user clicks on the position of the button with a mouse, the display system cannot determine whether the click operation is for the button or for the map. This makes it difficult to determine the destination of the event indicative the click operation.\nA simple solution is preparing an event-only filter for each layer. However, this solution is not practical, because the producer of its application must set the filters in consideration of the arrangement of the layers and the objects.\nAnother solution is using alpha values such as in a layer window, one of the user interfaces provided by Microsoft Windows® operating systems. The layer window allows mouse messages to pass through regions of an alpha value of zero. However, with the layered window, even a region with a low alpha value always receives a process request unless the alpha value is zero. When a plurality of alpha values other than zero is set for one layer, some regions of low alpha values are not regions selected by the user, so that they should not receive a process request. On the other hand, regions of high alpha values thus having clear images, such as text information, sometimes should not receive a process request."}
{"text": "The present invention relates generally to endoscopes and endoscopic procedures. More particularly, it relates to a method and apparatus to facilitate insertion of a flexible endoscope along a tortuous path, such as for colonoscopic examination and treatment.\nAn endoscope is a medical instrument for visualizing the interior of a patient\"\"s body. Endoscopes can be used for a variety of different diagnostic and interventional procedures, including colonoscopy, bronchoscopy, thoracoscopy, laparoscopy and video endoscopy.\nColonoscopy is a medical procedure in which a flexible endoscope, or colonoscope, is inserted into a patient\"\"s colon for diagnostic examination and/or surgical treatment of the colon. A standard colonoscope is typically 135-185 cm in length and 12-19 mm in diameter, and includes a fiberoptic imaging bundle or a miniature camera located at the instrument\"\"s tip, illumination fibers, one or two instrument channels that may also be used for insufflation or irrigation, air and water channels, and vacuum channels. The colonoscope is inserted via the patient\"\"s anus and is advanced through the colon, allowing direct visual examination of the colon, the ileocecal valve and portions of the terminal ileum.\nInsertion of the colonoscope is complicated by the fact that the colon represents a tortuous and convoluted path. Considerable manipulation of the colonoscope is often necessary to advance the colonoscope through the colon, making the procedure more difficult and time consuming and adding to the potential for complications, such as intestinal perforation. Steerable colonoscopes have been devised to facilitate selection of the correct path though the curves of the colon. However, as the colonoscope is inserted farther and farther into the colon, it becomes more difficult to advance the colonoscope along the selected path. At each turn, the wall of the colon must maintain the curve in the colonoscope. The colonoscope rubs against the mucosal surface of the colon along the outside of each turn. Friction and slack in the colonoscope build up at each turn, making it more and more difficult to advance, withdraw, and loop the colonoscope. In addition, the force against the wall of the colon increases with the buildup of friction. In cases of extreme tortuosity, it may become impossible to advance the colonoscope all of the way through the colon.\nSteerable endoscopes, catheters and insertion devices for medical examination or treatment of internal body structures are described in the following U.S. patents, the disclosures of which are hereby incorporated by reference in their entirety: U.S. Pat. Nos. 4,543,090; 4,753,223; 5,337,732; 5,337,733; 5,383,852; 5,487,757; 5,624,381; 5,662,587; and 5,759,151.\nAccordingly, an improved endoscopic apparatus is disclosed herein for the examination of a patient\"\"s colon or other internal bodily cavities with minimal impingement upon bodily cavities or upon the walls of the organs. A steerable endoscope having an elongate body with a manually or selectively steerable distal portion, an automatically controlled portion, which may be optionally omitted from the device, a flexible and passively manipulated proximal portion, and an externally controlled and manipulatable tracking rod or guide is described below. The tracking rod or guide may be slidably positioned within a guide channel or lumen within the endoscope or it may be externally positionable such that the guide and the endoscope may slide relative to one another along a rail or channel located along an external surface of the endoscope.\nIn operation, the steerable distal portion of the endoscope may be first advanced into a patient\"\"s rectum via the anus. The endoscope may be simply advanced, either manually or automatically by a motor, until the first curvature is reached. At this point, the steerable distal portion may be actively controlled by the physician or surgeon to attain an optimal curvature or shape for advancement of the endoscope. The optimal curvature or shape is considered to be the path which presents the least amount of contact or interference from the walls of the colon. In one variation, once the desired curvature has been determined, the endoscope may be advanced further into the colon such that the automatically controlled segments of controllable portion follow the distal portion while transmitting the optimal curvature or shape proximally down the remaining segments of the controllable portion. The operation of the controllable segments will be described in further detail below.\nAlternatively, once the steerable distal portion has been steered or positioned for advancement, the guide may be advanced distally in its flexible state along or within the endoscope until it reaches a distal position, i.e., preferably some point distal of the flexible proximal portion. Regardless whether the optional controllable portion is omitted or not from the device, the guide may be advanced near or to the end of the distal portion. Once the guide has been advanced, it may directly attain and conform to the curvature or shape defined by the steerable distal portion.\nPreferably, the guide is advanced to the distal end of steerable distal portion or, if the controllable portion is included in the device, the guide may be advanced to the distal end of the controllable portion, or to some point between the two portions. The guide may be advanced to any distal position as long as a portion of the guide attains and conforms to the optimal curvature or shape. Prior to advancing the endoscope over the guide, the guide may be left in its flexible state or it may be optionally rigidized, as discussed further below. If left in its flexible state, the guide may possibly provide desirable column strength to the endoscope as it is advanced through the colon over the guide. It is preferable, however, that the guide is rigidized once it has attained and conformed to the curvature. This allows the flexible proximal portion, i.e., the passive portion, to remain flexible and lightweight in structure. As the position of the guide is preferably rigidized and maintained, the endoscope may then be advanced over the guide in a monorail or xe2x80x9cpiggy-backxe2x80x9d fashion so that the flexible proximal portion follows the curve held by the guide until the endoscope reaches the next point of curvature.\nThis process of alternately advancing the guide and the endoscope may be repeated to advance the entire endoscope through the colon while the guide may be alternatively rigidized and relaxed while being advanced distally. While the endoscope is advanced through the colon, the physician or surgeon may stop the advancement to examine various areas along the colon wall using, e.g., an imaging bundle located at the distal end of the endoscope. During such examinations, the guide may be temporarily withdrawn from the endoscope to allow for the insertion of other tools through the guide channel if there is no separate channel defined within the endoscope for the guide. The guide may also be withdrawn through the instrument to any location within the body of the endoscope. In other words, the guide may be withdrawn partially or removed entirely from the endoscope at any time, if desired, because there are no constraints which may limit the travel of the guide through the body of the endoscope. After a procedure has been completed on the colon wall, the tool may be withdrawn from the guide channel and the guide may be reintroduced into the endoscope so that the endoscope may optionally be advanced once again into the colon.\nA further variation on advancing the endoscope may use multiple guides which are alternately rigidized while being advanced distally along a path. Although multiple guides may be used, two guides are preferably utilized. As the endoscopic device approaches a curvature, a first guide may be advanced in a relaxed and flexible state towards the steerable distal end of the device. While being advanced, the first guide preferably conforms to the shape defined by the distal end and the first guide may be subsequently rigidized to maintain this shape. The device may then be advanced further distally along the pathway while riding over the rigidized first guide.\nAfter the device has been advanced to its new position, a second guide may also be advanced distally in its relaxed state through the device up to the distal end while the first guide is maintained in its rigidized state. The second guide may then conform to the new shape defined by the distal end of the device and become rigidized to maintain this new shape. At this point, the first guide is also preferably maintained in its rigid state until the distal end of the device has been advanced further distally. The first guide may then be relaxed and advanced while the rigidity of the second guide provides the strength for advancing the guide. This procedure may be repeated as necessary for negotiating the pathway.\nTo withdraw the endoscope from within the colon, the procedure above may be reversed such that the withdrawal minimally contacts the walls of the colon. Alternatively, the guide may simply be removed from the endoscope while leaving the endoscope within the colon. Alternatively, the guide may be left inside the endoscope in the relaxed mode. The endoscope may then be simply withdrawn by pulling the proximal portion to remove the device. This method may rub or contact the endoscope upon the walls of the colon, but any impingement would be minimal.\nThe selectively steerable distal portion can be selectively steered or bent up to a full 180xc2x0 bend in any direction. A fiberoptic imaging bundle and one or more illumination fibers may extend through the body from the proximal portion to the distal portion. The illumination fibers are preferably in communication with a light source, i.e., conventional light sources, which may be positioned at some external location, or other sources such as LEDs. Alternatively, the endoscope may be configured as a video endoscope with a miniaturized video camera, such as a CCD camera, positioned at the distal portion of the endoscope body. The video camera may be used in combination with the illumination fibers. Optionally, the body of the endoscope may also include one or two access lumens that may optionally be used for insufflation or irrigation, air and water channels, and vacuum channels, etc. Generally, the body of the endoscope is highly flexible so that it is able to bend around small diameter curves without buckling or kinking while maintaining the various channels intact. The endoscope can be made in a variety of other sizes and configurations for other medical and industrial applications.\nThe optional controllable portion is composed of at least one segment and preferably several segments which may be controllable via a computer and/or controller located at a distance from the endoscope. Each of the segments preferably have an actuator mechanically connecting adjacent segments to allow for the controlled motion of the segments in space. The actuators driving the segments may include a variety of different types of mechanisms, e.g., pneumatic, vacuum, hydraulic, electromechanical motors, drive shafts, etc. If a mechanism such as a flexible drive shaft were utilized, the power for actuating the segments would preferably be developed by a generator located at a distance from the segments, i.e., outside of a patient during use, and in electrical and mechanical communication with the drive shaft. A proximal portion comprises the rest of the endoscope and preferably a majority of the overall length of the device. The proximal portion is preferably a flexible tubing member which may conform to an infinite variety of shapes. It may also be covered by a polymeric covering optionally extendable over the controllable portion and the steerable distal portion as well to provide a smooth transition between the controllable segments and the flexible tubing of the proximal portion. The controllable portion may be optionally omitted from the endoscope. A more detailed description on the construction and operation of the segments may be found in U.S. patent application Ser. No. 09/969,927 entitled xe2x80x9cSteerable Segmented Endoscope and Method of Insertionxe2x80x9d filed Oct. 2, 2001, which has been incorporated by reference in its entirety.\nA proximal handle may be attached to the proximal end of the proximal portion and may include imaging devices connected to the fiberoptic imaging bundle for direct viewing and/or for connection to a video camera or a recording device. The handle may be connected to other devices, e.g., illumination sources and one or several luer lock fittings for connection to various instrument channels. The handle may also be connected to a steering control mechanism for controlling the steerable distal portion. The handle may optionally have the steering control mechanism integrated directly into the handle, e.g., in the form of a joystick, conventional disk controller using dials or wheels, etc. An axial motion transducer may also be provided for measuring the axial motion, i.e., the depth change, of the endoscope body as it is advanced and withdrawn. The axial motion transducer can be made in many possible configurations. As the body of the endoscope slides through the transducer, it may produce a signal indicative of the axial position of the endoscope body with respect to the fixed point of reference. The transducer may use various methods for measuring the axial position of the endoscope body.\nThe guide is generally used to impart a desired curvature initially defined by the steerable portion and/or by the optional controllable portion to the passive proximal portion when the endoscope is advanced. If advanced into the steerable portion, the guide is preferably advanced to or near the distal tip of the portion. It is also used to impart some column strength to the proximal portion in order to maintain its shape and to prevent any buckling when axially loaded. Preferably, the guide is slidably disposed within the length of the endoscope body and may freely slide entirely through the passive proximal portion, through the controllable portion, and the steerable distal portion. The extent to which the guide may traverse through the endoscope body may be varied and adjusted according to the application, as described above. Furthermore, the proximal end of the guide may be routed through a separate channel to a guide controller which may be used to control the advancement and/or withdrawal of the guide and which may also be used to selectively control the rigidity of the guide as controlled by the physician.\nThe structure of the guide may be varied according to the desired application. The following descriptions of the guide are presented as possible variations and are not intended to be limiting in their structure. For instance, the guide may be comprised of two coaxially positioned tubes separated by a gap. Once the guide has been placed and has assumed the desirable shape or curve, a vacuum force may be applied to draw out the air within the gap, thereby radially deforming one or both tubes such that they come into contact with one another and lock their relative positions.\nAnother variation on the guide is one which is rigidizable by a tensioning member. Such a guide may be comprised of a series of individual segments which are rotatably interlocked with one another in series. Each segment may further define a common channel through which a tensioning member may be positioned while being held between a proximal and a distal segment. During use, the tensioning member may be slackened or loosened enough such that the guide becomes flexible enough to assume a shape or curve defined by the endoscope. When the guide is desirably situated and has assumed a desired shape, the tensioning member may then be tensioned, thereby drawing each segment tightly against one another to hold the desired shape.\nAnother variation may use a guide which is comprised of interlocking ball-and-socket type joints which are gasketed at their interfaces. Such a design may utilize a vacuum pump to selectively tighten and relax the individual segments against one another. Other variations may include alternating cupped segments and ball segments, a series of collinear sleeve-hemisphere segments, as well as other designs which may interfit with one another in series. Such a guide may be tightened and relaxed either by tensioning members or vacuum forces.\nA further variation on the guide is a coaxially aligned stiffening member. This assembly may include a first subassembly comprising a number of collinearly nested segments which may be held by a tensioning member passing through each segment. The first subassembly may be rigidized from a flexible or flaccid state by pulling on this tensioning member. A second subassembly may comprise a number of annular segments also collinearly held relative to one another with one or more tensioning members passing through each annular segment. The second subassembly preferably defines a central area in which the first nested subassembly may be situated coaxially within the second subassembly. The first subassembly is preferably slidably disposed relative to the second subassembly thereby allowing each subassembly to be alternately advanced in a flexible state and alternately rigidized to allow the other subassembly to be advanced. This design presents a small cross-section relative to the endoscope or device through which it may be advanced."}
{"text": "1. Technical Field of the Invention\nThe present invention relates generally to digital image processing, and particularly to demosaicing methods for interpolating colors at each pixel location.\n2. Description of Related Art\nDigital color image sensors are predominately of two types: CCDs (Charge Coupled Devices) and CMOS—APS (Complimentary Metal Oxide Semiconductor—Active Pixel Sensors). Both types of sensors typically contain an array of photo-detectors (pixels), arranged in rows and columns, that sample color within an image. Each pixel measures the intensity of light within one or more ranges of wavelengths, corresponding to one or more colors.\nIn addition, both types of sensors may include a color filter array (CFA), such as the CFA described in U.S. Pat. No. 3,971,065 to Bayer (hereinafter referred to as Bayer), which is hereby incorporated by reference. With the Bayer CFA, each pixel sees only one color: red, green or blue. The sampling pattern is shown below.\n                    G                    R                    G                    R                    G                            B                    G                    B                    G                    B                            G                    R                    G                    R                    G               \nA sensor fitted with a Bayer CFA produces a mosaiced image that is sampled in both the color space and in the spatial domain. The sampling process produces aliasing artifacts in both the color space and in the spatial domain. For example, since the full color spectrum is sampled only at certain pixel locations (depending on the CFA colors), it is impossible to accurately reconstruct the true color of an image, thereby producing color space aliasing artifacts. In addition, since high spatial frequencies in the original image are sampled at too low of a frequency, the original high frequencies in the image cannot be restored later on through image processing, thereby producing spatial domain aliasing artifacts.\nOne solution to the color space and spatial domain aliasing artifact problems is demosaicing. Demosaicing is the process of interpolating colors of a digital image obtained from an image sensor to obtain all three primary colors at a single pixel location. When viewed, the resulting demosaiced image provides a better quality and more visually pleasing image than the mosaiced image.\nPrevious demosaicing methods have included, for example, pixel replication, bilinear interpolation and median interpolation. In pixel replication, each missing value is taken from the neighbor to the left, above, or diagonally above and left, whichever is nearest. Bilinear interpolation offers some improvement over pixel replication with a moderate increase in complexity. In the bilinear interpolation method, each missing value is calculated based on an average of the neighboring pixel values, horizontally, vertically and/or diagonally. Median interpolation, which is a nonlinear interpolation method, offers the best results among these three algorithms (pixel replication, bilinear and median), especially when there are defective pixels, but has the maximum complexity. Median interpolation has two steps. First, missing values having four diagonal neighbors are interpolated using the median of those four values. Second, the remaining missing pixels are interpolated by the median of north, south, east, and west neighbors.\nHowever, with all of the above demosaicing methods, noticeable aliasing artifacts appear at edges within the image, producing a “zipper” effect in the demosaiced image. The “zipper” effect makes a straight edge in the image appear as a zipper. Solutions to the “zipper” effect have been proposed in U.S. Pat. No. 5,652,621 to Adams, Jr. et al. (hereinafter referred to as the Adams-Hamilton algorithm) and Kuno et al., “Aliasing reduction method for color digital still cameras with a single-chip charge-coupled device,” Journal of Electronic Imaging, Vol. 8, No. 4, October 1999, pp. 457–466 (hereinafter referred to as the Kuno algorithm), both of which are hereby incorporated by reference.\nThe Adams-Hamilton interpolation algorithm seeks to optimize the performance of image systems that sample images having horizontal and vertical edges. In the Adams-Hamilton interpolation algorithm, an interpolation method for interpolating a missing color value is selected. The interpolation direction is selected by calculating horizontal and vertical (or positive diagonal and negative diagonal) classifier values that are based on values taken from the color plane being interpolated and the color plane of the pixel location.\nThe two classifier values are compared with predefined thresholds and each other to select an interpolation method. For example, if one or more of the classifier values are greater than a predefined threshold, the missing color value is interpolated from both neighboring values of the same color as the missing color value and neighboring values of the color of the pixel location. The missing color value may be interpolated from horizontal neighboring values, vertical neighboring values or both horizontal and vertical neighboring values depending on which classifier values exceed the predefined threshold. The Adams-Hamilton algorithm does preserve detail and reduce the appearance of the “zipper” effect. However, with the Adams-Hamilton algorithm, some “zippering” is still visible.\nThe Kuno interpolation algorithm is based on the correlation of details assumption, expressed by Kodera et al. in “Method to display full-color image by using the correlation of color,” IIEEJ Annual Convention, Vol. 16(2), pp. 86–88 (1988), which is hereby incorporated by reference. The correlation of details assumption presumes that in a local region of an image, there is a high correlation amongst the colors. For example, using an image sensor fitted with a Bayer CFA, the correlation of details assumption predicts that the red, blue and green color values within a local region of an image will be relatively similar (e.g., if the red value of a red pixel location is 10, the correlation of details assumption predicts that the neighboring green and blue pixel locations will produce respective green and blue color values near 10).\nThe Kuno algorithm further postulates that even in regions with sharp color changes, a color correlation exists in one direction (i.e., either horizontal or vertical). Therefore, the Kuno algorithm uses a vertical interpolation orientation where only a low correlation in color values exists in the horizontal direction, and a horizontal interpolation orientation where only a low correlation in color values exists in the vertical direction. Specifically, the Kuno algorithm computes the horizontal and vertical gradients of the color plane associated with the missing color value and compares the horizontal and vertical gradients to a predefined threshold to determine the direction of interpolation. For example, if both the horizontal and vertical gradients are less than the predefined threshold, interpolation is performed both horizontally and vertically. However, if either of the gradients is greater than the predefined threshold, the direction of interpolation is in the smaller of the two gradients.\nOnce the direction of interpolation is determined, the Kuno algorithm uses a division operation to interpolate the missing color value. As with the Adams-Hamilton algorithm, the Kuno algorithm does reduce the “zipper” effect and improve color resolution as compared with previous demosaicing methods, such as pixel replication, bilinear interpolation and median interpolation. However, with the Kuno algorithm, some “zippering” is still visible. In addition, the Kuno algorithm requires a division operation to be carried out at every pixel location, which is expensive to implement. Therefore, what is needed is a demosaicing algorithm that further reduces the visibility of the “zipper” effect, as compared with the Adams-Hamilton algorithm and Kuno algorithm, and that can be implemented without the hardware complexity necessitated by the Kuno algorithm."}
{"text": "The present invention relates to a kneading machine or a mixing machine in which two rotors are respectively housed in two cavities of a chamber.\nHeretofore, as an apparatus for kneading natural or synthetic rubber or for mixing and kneading the rubber and various kinds of additive agents to disperse them uniformly in the rubber, or for mixing and kneading various thermoplastic elastomer or high viscous fluid with powder material or liquid material, or for dispersing materials, has been known a kneading machine in which blade-like rotors are rotatably supported within a sealed container (chamber) and kneading objects put into the chamber are moved from the axial direction of the rotor to the radial direction of the rotor to be circulated. For example, in the sealed type kneading machine disclosed in JP-A-8-24833 (corresponding to U.S. Pat. No. 4,834,543), spiral blades having spiral angles extending in the same direction are respectively formed on circumferences of shafts of two rotors and the two rotors are respectively arranged in disengaged condition in two parallel cylindrical cavities communicated with each other at their side faces. The rotors are rotated in opposite directions and the kneading objects are moved in the axial direction in each of the cavities and are interchanged between the cavities.\nFIG. 7 is a side view showing the rotors described in the above described Patent Document, and FIGS. 8-10 are schematic views showing flows of the kneading objects in the kneading machine in which the rotors are incorporated. A pair of rotors 1 are horizontally arranged in the right and left cavities of the chamber 3 of the kneading machine with a predetermined distance between them while keeping parallel relationship with each other. Neck portions of the shafts of the rotors la are rotatably supported by both sides of the chamber 3, and end portions of the rotor shaft at one side are connected to a drive motor via a gear system (not shown). The rotors are rotated in the opposite directions at the same speed. These rotors 1 each are formed with a long blade 5 and a short blade 6 extending in the axial direction with a predetermined spiral angle. The rotors are assembled in the machine to shift phases of rotating angle so that the rotors cannot interfere with each other at rotation and in the example, the rotors are assembled in the machine with a phase difference 90.degree. at a position.\nIn this example, the spiral directions of the blades of the right and left rotors 1 are the same. Therefore, when the rotors 1 are rotated in opposite directions, the movements of the kneading objects in the axial direction in the respective cavities 2 become opposite direction as schematically illustrated in FIGS. 8 and 9, so that the movement of the kneading objects in the axial direction in the cavities becomes good condition. Further, interchanging of the kneading objects are performed near the end of both cavities 2 in the axial direction.\nKneading machines of the above described type in which spiral directions of the blades portions of the right and left rotors are the same and the rotors are rotated in opposite directions can improve the movability of the kneading objects in the axial direction in the chamber. However, when an operation in which the raw materials of kneading objects are continuously supplied into the chamber and simultaneously the product is simultaneously taken out from a discharge outlet of the chamber, so-called continuous operation, is performed, resulted is a disadvantageous phenomenon, so-called short pass, in which the raw materials entered a material supply inlet of the chamber immediately move to the discharge outlet without being uniformly mixed in the radial direction (see FIG. 10) because of the high movability in the axial direction of raw materials, so that uniformly kneaded product cannot be obtained. Therefore, the conventional kneading machines can do only batch treatment, and they have a problem on the productivity.\nFurther, in the conventional kneading machines, because the right and left rotors of which spiral directions of the blade portions the rotors are the same are rotated in opposite directions, phases of the right and left rotors during rotation are shifted at the positions in the axial direction. Therefore, quantity of interchanging of the kneading objects between the cavities in the chamber changes. In case where the length of the chamber in the axial direction is short, the phase shift is limited to a small range and its influence is negligible. However, in case of the length of the chamber in the axial direction is longer, i.e. the ratio L/D of the length L with the diameter D of the chamber is larger, its influence cannot be ignored."}
{"text": "Transducers convert signals from one domain to another and are often used in sensors. One common sensor with a transducer that is seen in everyday life is a microphone that converts sound waves to electrical signals.\nMicroelectromechanical system (MEMS) based sensors include a family of transducers produced using micromachining techniques. MEMS, such as a MEMS microphone, gather information from the environment by measuring the change of physical state in the transducer and transferring a transduced signal to processing electronics that are connected to the MEMS sensor. MEMS devices may be manufactured using micromachining fabrication techniques similar to those used for integrated circuits.\nMEMS devices may be designed to function as, for example, oscillators, resonators, accelerometers, gyroscopes, pressure sensors, microphones, and micro-mirrors. Many MEMS devices use capacitive sensing techniques for transducing the physical phenomenon into electrical signals. In such applications, the capacitance change in the sensor is converted to a voltage signal using interface circuits.\nOne such capacitive sensing device is the MEMS microphone. A MEMS microphone generally has a deflectable membrane separated by a small distance from a rigid backplate. In response to a sound pressure wave incident on the membrane, it deflects towards or away from the backplate, thereby changing the separation distance between the membrane and backplate. Generally, the membrane and backplate are made out of conductive materials and form “plates” of a capacitor. Thus, as the distance separating the membrane and backplate changes in response to the incident sound wave, the capacitance changes between the “plate” and an electrical signal is generated.\nGenerally, a MEMS microphone sends the electrical signals generated in response to incident sound waves to some type of interface and/or processing electronics that may perform different functions depending on the application. As applications advance and increase in number, numerous challenges present themselves that produce inventive solutions for improving and innovating transducer systems and interface electronics."}
{"text": "Wideband Code Division Multiple Access (WCDMA) telecommunication systems have many attractive properties that can be used for future development of telecommunication services. In particular, the enhanced uplink of the WCDMA system is one ingredient in the mobile broadband solution of WCDMA. Now, in order to retain stability of a WCDMA cell, possibly in a cell running enhanced uplink, the load needs to be kept below a certain level. This follows since the majority of uplink user channels, at least in WCDMA, are subject to power control. This power control aims at keeping the received power level of each channel at a certain signal to interference ratio (SIR), in order to be able to meet specific service requirements.\nSince the Radio Base Station (RBS) tries to keep each channel at its specific preferred SIR value, it may happen that an additional user, or bursty data traffic of an existing user, raises the interference level, thereby momentarily reducing the SIR for the other users. The response of the RBS is to command a power increase to all other users, something that increases the interference even more. Normally this process remains stable below a certain load level. In case a high capacity channel would suddenly appear, the raise in the interference becomes large and the risk for instability, a so called power rush, increases. It is thus a necessity to schedule high capacity uplink channels, like the enhanced uplink (EUL) channel in WCDMA, so that one can insure that instability is avoided. In order to do so, the momentary load must be estimated in the RBS or any node connected thereto. This enables the assessment of the capacity margin that is left to the instability point.\nTo handle increasing uplink data rates, interference suppression (IS) and interference cancellation (IC) are being introduced in WCDMA. The main result is a reduction in the interference experienced by a user. The IS and IC also affect the load measurement functionality of the WCDMA enhanced uplink (EUL).\nThe instantaneous uplink load without IS/IC is today preferably obtained by a dedicated algorithm which estimates the rise over thermal (RoT), from measurements of the received total wideband power (RTWP) [1], [2]. There is also an uplink (UL) receiver structure concept defined that is a combination of the G-rake+ receiver and traditional interference cancellation (SIC) [3].\nPrior art methods for estimating load without IC, RoT estimation algorithms without IS/IC and IC methods with regeneration and subtraction are discussed more in detail in Appendix A.\nThere are a number of IS/IC methods available in prior art. Frequency domain equalization (FDE) and Frequency domain pre-equalization (FDPE) are two examples, where the interference is handled in the frequency domain, i.e. they are frequency domain receiver techniques.\nMinimum mean square error (MMSE) FDE is a common technique to combat frequency selectivity in wideband channels, like in the WCDMA uplink. FDE can be seen as providing a counterpart to the G-rake receiver. The main advantages, as compared to time domain equalization, include a reduced complexity. FDE primarily provides (self-)interference suppression\nFDPE is another frequency domain technique for interference suppression that has been developed as an alternative to the G-rake+ receiver structure. The main advantages associated with FDPE as compared to G-rake+ are claimed to be a lower complexity, a simpler receiver structure where much can be reduced and algorithms that may re-use LTE ASIC accelerators for fast Fourier transforms. At the same time, most of the interference suppression gains associated with the G-rake+ receiver remain.\nAn overview of these techniques is given in Appendix B.\nThe prior art cell stability load estimation functionality, exploits load factors for each user. In their simplest form the load factors are given by:\n            L      u        =                            P          u                RTWP            =                                    (                          C              /              I                        )                    u                          1          +                                    (                              C                /                I                            )                        u                                ,      u    =    1    ,  …  ⁢          ,  U  ,where Pu is the power of user u. Load factors are then summed up, for each power controlled user. In this way the neighbor cell interference is not included in the resulting load measure. This is reasonable since the neighbor cell interference should not affect the own cell power control loop, at least not when first order effects are considered.\nHowever, using prior art solutions of [1], [2] for load estimation, the scheduler will experience the load before IS or IC, e.g. according to FDE or FDPE, is applied. Hence the scheduling performance will then be the same as without IS or IC, i.e. under-scheduling and under-utilization will result.\nA problem with existing load estimation algorithms compatible with the FDE and FDPE receiver structures, known in public prior art, is that they overestimate the air-interface load of the uplink, relevant for cell stability. Thereby they cause under-scheduling in the RBS, which results in a too low throughput and/or capacity. They also cause blocking in the admission control function in the RNC, which also results in reduced throughput/capacity. Another problem is also that the admission and congestion control algorithms that reside in the RNC are not able to admit users so that the uplink with IS/IC is fully exploited."}
{"text": "Bacillus thuringiensis, Berliner is well-known for use as a microbial insect pathogen useful against leaf-chewing insects such as the alfalfa caterpiller, the cosmopolitan green beetle, the European corn borer and Mediterranean flower moth. See, for example, U.S. Pat. 3,150,062. Although used for the toxic effects on general agricultural and forest leaf-chewing insect pests, Bacillus thuringiensis, Berliner has no substantial effect on sucking insects such as aphids and the like.\nThe search has continued for ways to increase the potency, toxicity and/or persistance of insecticides including Bacillus thuringiensis, Berliner. It has been proposed, for example, in Chemical Abstracts, 28409v, Vol. 79, p. 148 (1973) to mix Bacillus thuringiensis with equal amounts of a synthetic pyrethroid insecticide. Such mixtures, however, have been found undesirable since, for example, while Bacillus thuringiensis is exempt from tolerance levels, the synthetic pyrethroid insecticide is not and the resulting mixture containing a synthetic pyrethroid insecticide also is not exempt.\nThe search has continued for a more efficacious manner of utilizing Bacillus thuringiensis, Berliner as a general insecticide."}
{"text": "A proline acylase (N-acyl-L-proline amido hydrolase) from a pseudomonas strain is described in Biochimica et Biophysica Acta, 744 (1983), pp. 180-188, Elsevier Biomedical Press. That enzyme is very rapidly inactivated at a temperature as low as 50.degree. C. However, enzymes are required for industrial use which have a high degree of stability. It is especially advantageous to carry out the reaction at elevated temperatures, which considerably raise the reaction speed and the solubility of the substrates, but this requires an enzyme which can withstand such temperatures."}
{"text": "Solid state power amplifiers are advantageous for their compact size and easy integration into semiconductor circuit components. Unfortunately, methods of manufacture for present day semiconductor power amplifiers require a semiconductor substrate dedicated to power amplifier devices or many processing steps in addition to common semiconductor processing steps for typical semiconductor complementary metal-oxide-semiconductor (CMOS) devices or their variants.\nFor example, high-end power amplifiers are built in gallium arsenide (GaAs) technologies, which require a GaAs substrate and dedicated processing steps that are not compatible with silicon-based CMOS technologies. As a result, the power amplifiers that utilize GaAs technologies tend to be costly. Middle-range power amplifiers are built in modified silicon germanium bipolar complementary metal-oxide-semiconductor (SiGe BiCMOS) technologies developed for high voltage power applications. Even modified SiGe BiCMOS technologies tend to add its own cost associated with enabling power amplifiers. Enabling power amplifiers in standard CMOS technologies also tends to introduce many new processing steps and device modifications to accommodate the high voltages that the power amplifiers require, thus also increasing the manufacturing cost for the power amplifiers.\nA junction field effect transistor (JFET) is a semiconductor device in which the current between a source and a drain is controlled by the voltage applied to a junction gate terminal, or a “gate.” Unlike a metal-oxide-semiconductor field effect transistor (MOSFET), the gate of a JFET is not insulated from the source and the drain. Instead, the body of the transistor and the gate of the transistor form a reverse-biased pn junction with depletion regions both in the gate and in the body. Therefore, the JFET is a depletion mode device with a high input impedance. The input signal is supplied to the gate, typically in the form of a voltage input. The output is the current between the source and the drain which is modulated by the input voltage at the gate.\nA typical JFET comprises a source and a drain that are heavily doped with dopants of a first conductivity type, i.e., p-type or n-type, at a peak dopant concentration typically in the range from about 1.0×1020/cm3 to about 3.0×1021/cm3. The body of the JEFT is also doped with dopants of the first conductivity type at a dopant concentration typically in the range from about 1.0×1017/cm3 to about 1.0×1019/cm3. A channel is formed within the body along the pn junction boundary. The gate, located on the body and separated from the source and the drain, is heavily doped with dopants of a second conductivity type, which is the opposite type of the first conductivity type, at a peak dopant concentration typically in the range from about 1.0×1020/cm3 to about 3.0×1021/cm3. A voltage bias is applied across a gate contact and a body contact to form a reverse biased pn junction between the gate and the body. The gate contact and the body contact directly contact the gate and the body, respectively, and are typically a metal semiconductor alloy.\nOn a circuit level, the JFET gate presents a small current load, which is the reverse bias leakage of the gate-to-channel junction. The current load of a JFET, i.e., the gate current, is higher than the current load of a typical MOSFET, since the MOSFET has an extremely low gate current, for example, in the range of picoamperes, due to an insulator between the gate and the channel, i.e., a gate dielectric. However, the gate current of a typical JFET is much lower compared to the base current of a typical bipolar junction transistor (BJT), and the transconductance of a typical JFET is higher than that of a typical MOSFET, enabling handling of a higher current. For this reason, JFETs are used in high-input impedance linear amplifier circuits. Use of JFETs as a switch in power semiconductor circuits is also known.\nA high on/off impedance ratio is necessary in a JFET to enable high power amplification. To provide such a high on/off impedance ratio, a JFET needs to have a low impedance during the on state, while having a high impedance during the off state. To decrease the impedance in the on state, the cross-sectional area of the channel needs to be increased in a JFET. At the same time, leakage current through the channel needs to be minimized to increase the impedance in the off state.\nTherefore, there exists a need for a semiconductor structure that provides high power amplification and has a high on/off impedance ratio, and methods of manufacturing the same.\nSpecifically, there exists a need for a JFET structure having a high on current for power amplification as well as a high impedance in an off state, and methods of manufacturing the same."}
{"text": "Various technologies exist for controlling the appearance and relative position/scale (layout) of multiple media sources on a display. For example, video wall scalars and switchers allow a user to inject control commands into a video routing system that impacts the appearance of the video sources on a single shared display. Video wall control commands necessarily modify the shared appearance of the main display—that is, if a media source is scaled larger on the display—it is scaled for all viewers of that single shared display. At the other end of the spectrum, some systems support individual control that is different from the shared view. For example, some video teleconferencing software systems allow individual users to scale a video stream of a document or web camera view in their own interface, in the absence of a shared view.\nAs a result, these existing technologies do not address the needs of a distributed visualization system that includes both a host display system that provides a synchronous view of the various media streams and individual (but synchronized) views of that same system on each of the connected source computers. What is needed is a mechanism that allows users to control a shared display surface and to simultaneously view that shared display on connected devices in a way the supports group collaboration and individual viewing of each of the shared media sources, without modifying the main shared display screen.\nConsider the case, for example, when four source computers, each connected to a shared display a sharing media sources simultaneously. Assume the shared display is a 1920×1080 resolution screen. If each computer is sharing a single 1920×1080 resolution media stream and all four sources are being shown on the shared display equally, then, at most, any individual source resolution is (1920/2)×(1080/2)=960×540. While this is useful for comparative viewing of different, typically disparate, sources of media, information is necessarily lost. Furthermore, systems that display a replica of the shared display media on each of the individual computers may further reduce the resolution of each displayed media source so that it fits within a constrained user interface."}
{"text": "LCSs (sometimes referred to as load control receivers (LCRs)) are used as part of a demand-side management scheme to control the operation of certain power-consuming devices. Examples of such power-consuming devices include, but are not limited to, an air conditioning unit, a water heater, a heater, and a circulation fan. Controlling such power-consuming devices involves determining when a control event (i.e., preventing power from flowing to such devices) occurs and for how long.\nAt times, a LCS is installed with an existing power-consuming device. In such a case, the LCS is installed by electrically coupling different components of the existing power-consuming device to various portions of the LCR. For example, a power source from the power-consuming device is electrically coupled to the LCR. Such work often requires a licensed electrician because at least a portion of the power feeding the LCS is at least 110 V alternating current (AC), and so applicable regulations often require a licensed electrician when working with wiring at such voltages. Consequently, the installation of LCSs can be expensive and inconvenient."}
{"text": "Pruritus is defined as an unpleasant sensation that triggers an itch, namely a desire to scratch. Itch, like pain, is one of the body's basic defense mechanisms. A fundamental biologic function of itch is to alert an animal to the presence of potentially harmful toxins or other hazards such as disease-carrying insects, and to stimulate a reflex aimed at getting rid of these hazards. Itch can manifest acutely, like the reflex to remove fleas and other parasites. Alternatively, chronic itch, like pain, can be become self-perpetuating and pathologic in itself. Chronic itch results when peripheral and central nerves are over-stimulated, which leads to activation and proliferation of pruritus-mediating nerve fibers. Sensitized nerve fibers have been shown to more readily stimulate pruritus. Chronic itch necessitates more than symptomatic treatment, requiring a thorough diagnostic work-up to identify the underlying cause, and multimodal therapy to manage the insidious effects.\nProtein kinases are families of enzymes that catalyze the phosphorylation of specific residues in proteins, broadly classified into tyrosine and serine/threonine kinases. Inappropriate kinase activity, arising from mutation, over-expression, or inappropriate regulation, dysregulation, or de-regulation, as well as over- or under-production of growth factors or cytokines has been implicated in many diseases, including but not limited to cancer, cardiovascular diseases, allergies, asthma and other respiratory diseases, autoimmune diseases, inflammatory diseases, bone diseases, metabolic disorders, and neurological and neurodegenerative disorders such as Alzheimer's disease. Inappropriate kinase activity triggers a variety of biological cellular responses relating to cell growth, cell differentiation, survival, apoptosis, mitogenesis, cell cycle control, and cell mobility implicated in the aforementioned and related diseases. Thus, protein kinases have emerged as an important class of enzymes as targets for therapeutic intervention. In particular, the JAK family of cellular protein tyrosine kinases (JAK-1, JAK-2, JAK-3, and Tyk-2) play a central role in cytokine signaling (Kisseleva et al, Gene, 2002, 285, 1; Yamaoka et al. Genome Biology 2004, 5, 253)). Upon binding to their receptors, cytokines activate JAK which then phosphorylate the cytokine receptor, thereby creating docking sites for signaling molecules, notably, members of the signal transducer and activator of transcription (STAT) family that ultimately lead to gene expression, which stimulates biologic responses such as an itch signal. Activation of the JAK-STAT pathway also results in several other ancillary biologic activities that contribute to the inflammation and pruritic processes that contribute to acute allergy in animals but can also exacerbate clinical signs and contribute to chronic allergy.\nThere are substantial needs for safe and efficacious agents to control atopic dermatitis in animals, including mammals, birds, and fish. Examples of mammals include, but are not limited to, humans, cattle, sheep, goats, llamas, alpacas, pigs, horses, donkeys, dogs, cats, and other livestock or domestic mammals. Examples of birds include turkeys, chickens, ostriches, and other livestock or domestic birds. The market for treating atopic dermatitis in animals is currently dominated by corticosteroids, which cause distressing and undesirable side effects in animals, specifically in companion animals such as dogs. Antihistamines are also used, but are poorly effective. A formulation of cyclosporine (ATOPICA™) is currently being marketed for atopic dermatitis in dogs and cats, but is expensive and has a slow onset of efficacy. In addition, there are GI toleration issues with ATOPICA™. WO 2010/020905 discloses JAK inhibitors for the treatment of pruritus."}
{"text": "(1) Field of the Invention\nThe present invention relates generally to underwater weapons systems, and more particularly to an underwater weapon system having a rotating gun system mounted in a vehicle housing.\n(2) Description of the Prior Art\nUnderwater gun systems are being developed for use as anti-mine and anti-torpedo applications. The basic systems include an underwater gun, underwater projectiles, a ship-mounted turret supporting the underwater gun, a targeting system and a combat system. The underwater gun shoots the underwater projectiles which are specially designed for neutralization of undersea targets at ranges of 200 meters or more. The undersea targets are identified and localized with the targeting system, and the combat system provides the control commands to direct the ship-mounted turret to aim the gun towards the target.\nCurrently, underwater gun systems are mounted onboard a conventional ship by means of a retractable gun turret that deploys a rotating gun using a large hydraulic actuator. The gun is stored in a bay having a hatch on the bottom of the ship that can be opened to allow the turret and gun to be deployed therethrough. When deployed, the gun turret is not readily faired thereby causing drag that results in substantial hydrodynamic loads on the deployed turret and the host ship. Furthermore, for the gun system to respond in a rapid fashion to an unexpected threat, the gun system must be continuously deployed thereby subjecting the host ship to long-term hydrodynamic loads while also reducing the value of having a retractable system. Still further, the space required onboard the host ship is considerable, e.g., the length required for a retractable turret can exceed ten meters for a gun that fires 30 millimeter caliber rounds."}
{"text": "The present invention relates to a control device that controls a drive device in which a first engagement device, a rotary electric machine, and a second engagement device are provided on a power transfer path connecting between an input member drivably coupled to an internal combustion engine and an output member drivably coupled to wheels and are arranged in this order from a side of the input member."}
{"text": "The quality of the mouldings depends largely on the filling ratio of an injection mould. The filling ratio is defined by the shot capacity and this in turn depends on the metered volume and the injection pressure. Disturbing factors such as differences in the consistency of the moulding material or temperature fluctuations in the plasticizing chamber of the moulding machine or in the mould affect the metered volume and thus the filling ratio and mould quality."}
{"text": "For a proper fit e.g. of a hearing device it has always been considered vital to obtain good data of the ear canal. In case of hearing instruments for hard of hearing patients the data needs to cover a significant part of the canal and preferably beyond the “second bend”. Traditionally, data has been obtained with a silicon ear impression and this process is well-established. Being a manual and quite difficult process, alternative technologies have been searched for a long time and currently the focus is on direct scanning of the ear. An intermediate stage is scanning of ear impressions for computer-based in-the-ear (ITE) shell or behind-the-ear (BTE) mold production.\nThe pursuit of direct ear scanning is still ongoing and there are several approaches being studied. A number of problems remain to be solved, one being the way to insert the scanning head (probe) into the ear canal without hurting the patient, changing the shape of the canal or impacting the data acquisition. This is one of the subjects of the present invention.\nIn this respect in WO 02/091920 A1 the insertion of a scanning probe into an ear canal is proposed. Therefore, this prior art document is highly relevant and interesting, but does not give any contribution to the above mentioned problem.\nIn WO 02/16867 A1 a 360° probe-shaped non-contact scanning device is described. It also includes the post-processing of the data to obtain high-precision 3D data for individual ear canal shapes. But again nothing in this prior art document is addressing the issue relating to bends within the ear canal and insertion depth control.\nIn WO 02/16865 A2 a calibration method for ear scanning probes in order to achieve needed precision is described. The invention itself again does not touch the topic of the present invention.\nThere are several still open issues and real problems pursuing replacement of traditional ear impression taking by “direct ear scanning” one way or another. Most approaches involve the use of light (LED, laser, etc.); some propose video, ultrasound, computer tomography and even MRI (magnetic resonance imaging). Challenges relate to issues like: resolution, physical reference points, ear canal skin surface condition, post-processing of 3D information and touchless insertion/extraction of probes. The latter problem has not really been solved by any significant patent documents or publications to date and remains an obstacle in obtaining good 3D data of the undistorted ear canal. In addition, there are problems like skin scratches during insertion and extraction and even more severe: harm to the tympanic membrane, which is both very unpleasant and bears the risk of a distorted tympanic membrane function (e.g. conductive hearing loss).\nBesides the described problem in relation particularly to insertion into the ear canal, this of course is a general problem in relation to the insertion of any kind of scanning probes, sensors, video heads, tubes like cannulas or catheters etc. into orifices, canals, tubes and the like of the human or animal body. Besides the ear canal this could be e.g. the intestinal tract, the gullet, etc."}
{"text": "WO 2009/098623 A1 discloses a magnetic biosensor based on magnetic beads that can be actuated with electromagnetic fields. The magnetic beads are functionalized with antibodies that can bind a specific analyte molecule in a sample. The beads are attracted to the sensor surface, where the number of bound beads is directly or inversely related to the amount of analyte molecules present in the sample. The beads are then detected by a technique which is based on frustrated total internal reflection (FTIR)."}
{"text": "(1) Field of the Invention\nThe present invention relates to a mobile installation for loading, transporting and unloading such track parts as ties on, in and from open top railroad cars, which may be used in combination with a mobile apparatus for sequentially exchanging groups of old ties of a railroad track supporting the mobile apparatus, for example groups of one to three old ties between groups of one to three retained ties, for groups of new ties. The mobile installation comprises a train mounted for mobility along the railroad track and includes a plurality of the open top railroad cars, adjacent ones of the railroad cars being coupled together and each railroad car having two parallel side walls with top edges and two end walls, the end walls of the adjacent railroad cars defining respective gaps therebetween, a power-driven crane for loading and unloading the track parts and having an undercarriage supporting the crane for mobility in the direction of the railroad track, and a track supporting the undercarriage of the crane for mobility above the top edges of the railroad cars in this direction, the track comprising two parallel guide rails mounted on the top edges of the railroad cars and being spaced apart a distance corresponding to the gage of the crane undercarriage, the guide rails extending beyond the end walls of the railroad cars into the gaps between the adjacent cars.\n(2) Description of the Prior Art\nU.S. Pat. Nos. 4,175,902, dated Nov. 27, 1979, and 4,190,394, dated Feb. 26, 1980, disclose an apparatus and method for loading and unloading open top or gondola railroad cars used in the exchange of consecutive groups of ties. The apparatus comprises a train mounted for mobility along the railroad track and includes a plurality of the open top railroad cars having a considerable loading volume, adjacent ones of the railroad cars being coupled together and each railroad car having two high parallel side walls with top edges and two high end walls, the end walls of the adjacent railroad cars defining respective gaps therebetween, and a power-driven crane with booms for loading and unloading the ties and having two undercarriages supporting the crane for mobility in the direction of the railroad track. The undercarriages have pneumatic tires to enable the crane to be moved along a road or the railroad track, and the crane also has pivoted gliding feet or brackets for gripping the top edges, the relatively widely spaced top edges of the railroad cars serving as a track for moving the crane along the cars while the gliding feet grip the top edges. A cable is attached to the crane to pull the crane along the railroad cars as it is perched atop the cars. Operation of this apparatus requires great skill and a number of sometimes life-threatening manual steps. The movement of the crane between adjacent cars is particularly difficult and very time-consuming, which considerably reduces the efficiency of the operation. The crane movement along and between the cars is quite unstable, providing unsafe operating conditions and frequent interruptions. In addition, the tractor used for the crane must be specially designed to enable the crane to effectuate the required forward and rearward movements on the top edges of the gondola cars.\nU.S. Pat. Nos. 4,096,954, dated June 27, 1978, and 4,099,635, dated July 11, 1978, disclose a gondola car loader and tie handler, as also described on page 68 of \"Progressive Railroading\", February 1978, which also uses a power-driven crane movable atop a train of gondola cars. Two beams may be pivotally mounted on the chassis of the crane at each side thereof and these pivotal beams carry two flanged wheels for guidance of the crane chassis along the top edges of the gondola cars. The spacing between the guide wheels corresponds at least to the gap between adjacent cars. Each beam is pivotal at the center thereof for rotation about a transversely extending axis and is connected to a pivoting drive so that the front beams may be swung upwardly when the crane is moved from one to the next car, the crane being advanced until the wheels of the front beams may be engaged with the top edges of the next car. The beams are then swung down while the rear beams are pivoted upwardly and then lowered again when their wheels are engaged with the top edges of the next car. The operation of this apparatus is also slow and uneconomical, the movement of the crane being unstable and, therefore, unsafe for the operating personnel. Bridging the cars is not only difficult and time-consuming but is highly accident-prone in view of the heavy weight of the crane.\nThe above-described mobile installations with a power-driven tie loading and unloading crane movable atop a train of gondola cars have such a low efficiency that they cannot be used economically in tie exchange operations since they would cause long sections of track to be blocked for long periods for normal train traffic.\nU.S. Pat. No. 1,879,203, dated Sept. 27, 1932, discloses a train of open top railroad cars for loading, transporting and unloading bulk material. A self-propelled carriage for the bulk material runs on rails atop the cars, and a pivoted link arrangement is interposed between the rails in the gaps between adjacent cars to provide a continuous track from one end of the train to the other."}
{"text": "Measurement of a polarization mode dispersion of optical fibers has been carried out conventionally. For example, Japanese Patent Laid-Open Publication (Kokai) No. H09-264814 describes a polarization mode dispersion measuring device for optical fibers. With reference to FIG. 3, a description will now be given of the device for measuring the polarization mode dispersion of optical fibers according to Japanese Patent Laid-Open Publication (Kokai) No. H09-264814.\nFirst, the polarization mode dispersion τPMD of an optical fiber under test 104 is defined by the following equation (1):τPMD=2√{square root over ({dot over (θ)}2+{dot over (ψ)}12 cos2 θ+{dot over (ψ)}22 sin2 θ  (1)\nwhere θ is a polarization angle, ψ1 is a phase shift in a certain direction on a plane perpendicular to the propagation direction of light, and ψ2 is a phase shift in a direction orthogonal to ψ1. On this occasion, a transfer function matrix [T] of the optical fiber under test 104 is defined by the following equation (2):                               [                      T            ⁡                          (              ω              )                                ]                =                  [                                                                                                                                    T                      11                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        11                                                                                                                                                                                                        T                      12                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        12                                                                                                                                                                                                                                T                      21                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        21                                                                                                                                                                                                        T                      22                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        22                                                                                                                          ]                                    (        2        )            \nwhere |Tij| is an amplitude of respective matrix elements, φij is a phase shift of the respective matrix elements, and both of them are functions of an optical angular frequency ω. Then, the parameters θ, ψ1, and ψ2 in the equation (1) are respectively obtained by the following equations (3), (4), and (5).θ(ω)=0.5 cos−1(|T11|2−|T21|2)  (3)ψ1(ω)=(φ11−φ22)/2  (4)ψ2(ω)=(φ21−φ12+π)/2  (5)\nConsequently, the polarization mode dispersion τPMD of the optical fiber under test 104 is obtained by obtaining the transfer function matrix [T] of the optical fiber under test 104.\nA description will now be given of how to obtain the transfer function matrix [T] of the optical fiber under test 104 with reference to FIG. 3. First, a control unit 109 makes output light of a polarization controller 103 as a linearly polarized wave in line with a p direction of a polarization beam splitter 105 incident to the optical fiber under test 104. On this occasion, output light from the optical fiber under test 104 is represented by the following equation (6).                               [                                                                                                                                    T                      11                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        11                                                                                                                                                                                                        T                      12                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        12                                                                                                                                                                                                                                T                      21                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        21                                                                                                                                                                                                        T                      22                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        22                                                                                                                          ]                ⁢                                                       [                                                                    1                                                                                        0                                                              ]                        =                          [                                                                                                                                                                    T                          11                                                                                            ⁢                                              ⅇ                                                                              -                            j                                                    ⁢                                                                                                          ⁢                                                      ϕ                            11                                                                                                                                                                          (                                              p                        ⁢                                                                                                                                                                                                                                                                                                                            )                                                                                                                                                                                                            T                          21                                                                                            ⁢                                              ⅇ                                                                              -                            j                                                    ⁢                                                                                                          ⁢                                                      ϕ                            21                                                                                                                                                                          (                                              s                        ⁢                                                                                                                                                                                                                                                                                                                            )                                                                                  ]                                                          (        6        )            \nThe above-described output light is split into an s-polarization component and a p-polarization component by the polarization beam splitter 105, the components are made incident to O/E converters 1061 and 1062 respectively, and the O/E converters 1061 and 1062 respectively measure:|T11|e−jφ11, |T21|e−jφ21\nAfter the above-described measurement, the control unit 109 rotates the output light of the polarization controller 103 by 90°, and makes the light as a linearly polarized wave coincide with an s direction of the polarization beam splitter 105 incident to the optical fiber under test 104. On this occasion, output light from the optical fiber under test 104 is represented by the following equation (7).                               [                                                                                                                                    T                      11                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        11                                                                                                                                                                                                        T                      12                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        12                                                                                                                                                                                                                                T                      21                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        21                                                                                                                                                                                                        T                      22                                                                            ⁢                                      ⅇ                                                                  -                        j                                            ⁢                                                                                          ⁢                                              ϕ                        22                                                                                                                          ]                ⁢                                                       [                                                                    0                                                                                        1                                                              ]                        =                          [                                                                                                                                                                    T                          12                                                                                            ⁢                                              ⅇ                                                                              -                            j                                                    ⁢                                                                                                          ⁢                                                      ϕ                            12                                                                                                                                                                          (                                              p                        ⁢                                                                                                                                                                                                                                                                                                                            )                                                                                                                                                                                                            T                          22                                                                                            ⁢                                              ⅇ                                                                              -                            j                                                    ⁢                                                                                                          ⁢                                                      ϕ                            22                                                                                                                                                                          (                                              s                        ⁢                                                                                                                                                                                                                                                                                                                                                                                          )                                                                                  ]                                                          (        7        )            \nThe above-described output light is split into an s-polarization component and a p-polarization component by the polarization beam splitter 105, the components are made incident to the O/E converters 1061 and 1062 respectively, and the O/E converters 1061 and 1062 respectively measure:|T12|e−jφ12, |T22|e−jφ22\nA network analyzer 107 obtains the parameters θ, ψ1, and ψ2 from the respective parameters measured as described above, and the equations (3), (4), and (5). It should be noted that the network analyzer 107 controls an intensity modulation ratio in an optical intensity modulator 102 through an amplifier 108.\nThen, the above-described measurement is carried out while the output wavelength of a wavelength-variable light source 101 is being swept, thereby obtaining θ(ω), ψ1(ω), and ψ2(ω) from the respective measurement results. Then, the control unit 109 obtains the polarization mode dispersion τPMD from the equation (1).\nHowever, in the above-described method, it is necessary to cause the control unit 109 to switch the direction of the output light of the polarization controller 103 either to the p direction or the s direction of the polarization beam splitter 105. It is necessary to carry out the switching of the direction of the output light of the polarization controller 103 for each wavelength sweep of the wavelength-variable light source 101. Consequently, it takes a long period to measure the polarization mode dispersion τPMD.\nIn view of the foregoing, an object of the present invention is to provide a polarization mode dispersion measuring device and the like which reduce the time period required to measure the polarization mode dispersion τPMD."}
{"text": "The present invention is directed to the field of hydraulic dampers for reducing undesired pressure surges or oscillations in systems requiring hydraulic damping. This invention has particular applicability to the damping of periodically varying motions in systems which result in noise or oscillation in hydraulic pressure.\nIn a particular example, undesirable oscillation in a hydraulic system occurs in the nose wheel assembly of an aircraft, which is typically steered using a hydraulic system. The pivot point of the wheel is typically located at a point ahead of the wheel axle, rather like a caster. The mass of wheel is behind the pivot, and the dynamics of the system are such that the wheel tends to \"shimmy\" while moving along the ground, in a manner similar to other commonly-observed caster-type wheels, e.g. the wheels on a shopping cart.\nSuch wheel shimmy results in wear and tear on the wheel and supporting strut which may increase exponentially resulting structural damage to the aircraft. Suppression of shimmy often requires compromises in steering control accuracy. It is desirable to provide a hydraulic means of damping which does not compromise steering accuracy.\nOne known hydraulic damper uses a piston and a pair of springs to form a compliance, and a hydraulic orifice to form a damping control which reduces rapid shimmy motion while permitting the relatively slow motion of steering. However, such dampers include a number of mechanical elements.\nHydraulic systems are prone to pressure surges which occur when large valves are quickly opened or closed. This phenomenon, also known as \"water hammer,\" is also observed in common household plumbing when faucets are abruptly turned off, producing a sound and a shock to the pipes. Such surges also create additional wear and tear in hydraulic systems, thereby shortening the operational life of the respective components."}
{"text": "Pancreatitis is classified mainly into acute pancreatitis caused by autodigestion due to intracellular activation of the digestive enzymes produced in the exocrine pancreas and chronic pancreatitis which, by fibrosis, progresses to irreversible functional failure of the pancreas. The etiology of pancreatitis remains to be fully elucidated as yet. The therapeutic approach to acute pancreatitis consists of, to start with, conservative therapies such as removal of causes, pancreatic protection, arrest of pancreatic autodigestion, and removal of pain. Referring to pharmacotherapy, arrest of pancreatic autodigestion is sought by the administration of proteolytic enzyme inhibitors such as gabexate mesylate, camostat mesylate, nafamostat mesylate, urinastatin, and aprotinin. The treatment of chronic pancreatitis consists of therapies analogous to that of acute pancreatitis in relapsing episodes involving abdominal pain and prohibition of alcohol intake and restriction of fat-rich food in remissions. The protease inhibitors in use today for the pharmacotherapy of pancreatitis are only effective in relieving some of the symptoms of pancreatitis and have various side effects, thus being not clinically satisfactory enough.\nMeanwhile, the list of therapeutic drugs for pancreatitis under development today includes several compounds such as loxiglumide, which is a cholecystokinin antagonist, and FK480, which is a protease inhibitor. Furthermore, it has been reported that the antiulcer compound irsogladine maleate 2,4-diamino-6-(2,5-dichlorophenyl)-1,3,5-triazine maleate!, which is structurally akin to the compound of the invention, proved efficacious in an animal model of cerulein (hereinafter referred to briefly as Cn)-induced acute pancreatitis (the Proceedings of the 36th Congress of The Japanese Society of Gastroenterology, p.266).\nThe inventors of the present invention discovered that the 2-amino-4-substituted amino-6-(2,5-dichlorophenyl)-1,3,5-triazine derivative according to the invention has antihepatitis activity, thus being a compound of value as a therapeutic drug for hepatitis, and already filed a patent application (WO96/04914)."}
{"text": "Excavating equipment typically includes various wear parts to protect underlying products from premature wear. The wear part may simply function as a protector (e.g., a wear cap) or may have additional functions (e.g., an excavating tooth, which functions to break up the ground ahead of the bucket as well as protecting the underlying digging edge). In either case, it is desirable for the wear part to be securely held to the excavating equipment to prevent loss during use, and yet be capable of being removed and replaced when worn. In order to minimize equipment downtime, it is desirable for the worn wear part to be capable of being easily and quickly replaced in the field. Wear parts are usually formed of three (or more) components in an effort to minimize the amount of material that must be replaced on account of wearing. As a result, the wear part generally includes a support structure that is fixed to the excavating equipment, a wear member that mounts to the support structure, and a lock to hold the wear member to the support structure.\nAs one example, an excavating tooth includes an adapter as the support structure, a tooth point or tip as the wear member, and a lock or retainer to hold the point to the adapter. The adapter is fixed to the front digging edge of an excavating bucket and includes a nose that projects forward to define a mount for the point. The adapter may be a single unitary member or may be composed of a plurality of components assembled together. The point includes a front digging end and a rearwardly opening socket that receives the adapter nose. The lock is inserted into the assembly to releasably hold the point to the adapter.\nThe lock for an excavating tooth is typically an elongate pin member that is fit into an opening defined cooperatively by both the adapter and the point. The opening may be defined along the side of the adapter nose, as in U.S. Pat. No. 5,469,648, or through the nose, as in U.S. Pat. No. 5,068,986. In either case, the lock is inserted and removed by the use of a hammer. Such hammering of the lock can be an arduous task and impose a risk of harm to the operator.\nThe lock is usually tightly received in the passage in an effort to prevent ejection of the lock and the concomitant loss of the point during use. The tight fit may be effected by partially unaligned holes in the point and adapter that define the opening for the lock, the inclusion of a rubber member in the opening or in the pin, and/or close dimensioning between the lock and the opening. However, as can be appreciated, an increase in the tightness in which the lock is received in the opening exacerbates the difficulty and risk attendant with hammering the locks into and out of the assemblies.\nThe lock additionally often lacks the ability to provide substantial tightening of the point onto the adapter. While rubber members have been provided in prior locking systems to provide some tightening of the wear member on the support structure, it has tended to provide only limited benefit as the rubber lacks the strength needed to ensure a tight fit when the teeth are under load during use. Most locks also fail to provide any ability to be retightened as the parts become worn. As a result, many locks used in teeth are susceptible to being lost as the parts wear and the tightness decreases. Prior locks that provide take up or the ability to be retightened tend to rely upon threads or wedges, which commonly suffer from removal difficulties and/or safety issues.\nShortcomings in the locking arrangements are not limited strictly to the mounting of points on adapters. In another example, an adapter is a wear member that is fit onto a lip of an excavating bucket, which defines the support structure for the adapter. While the point experiences the most wear in the system, the adapter will also wear and in time need to be replaced. It is common for adapters to be mechanically attached to a bucket lip so as to permit the use of harder steel and to accommodate replacement in the field. One common approach is to use a Whisler style adapter, such as disclosed in U.S. Pat. No. 3,121,289 (see FIG. 8). In a traditional Whisler system, the adapter is formed with bifurcated legs that straddle the bucket lip. The adapter legs and the bucket lip are formed with openings that are aligned for receiving the lock. The lock in this environment comprises a generally C-shaped spool and a wedge. The arms of the spool overlie ramps on the rear end of the adapter legs. The ramps on the legs and the inner surfaces of the arms are each inclined rearward and away from the lip. The wedge is then hammered into the aligned openings to force the spool rearward. This rearward movement of the spool causes the arms to tightly pinch the adapter legs against the lip to prevent movement or release of the adapter during use.\nHowever, the hammering of the wedge into and out of the openings in a Whisler-style lock tends to be difficult and potentially hazardous. Removal can be particularly difficult as the bucket must generally be turned up to provide access for driving the wedges out of the assembly. In this orientation of the bucket the worker must access the opening from beneath the bucket and drive the wedge upward with a large hammer. The risk is particularly evident in connection with large buckets. Also, because wedges can eject during service, it is common for the wedges to be tack-welded to its accompanying spool, which eliminates any retightening and makes wedge removal more difficult.\nIn many assemblies, other factors can further increase the difficulty of removing and inserting the lock when replacement of the wear member is needed. For example, the closeness of adjacent components, such as in laterally inserted locks (see, e.g., U.S. Pat. No. 4,326,348), can create difficulties in hammering the lock into and out of the assembly. Fines can also become impacted in the openings receiving the locks making access to and removal of the locks difficult.\nThere have been some efforts to produce non-hammered locks for use in excavating equipment. For instance, U.S. Pat. Nos. 5,784,813 and 5,868,518 disclose screw driven wedge-type locks for securing points to adapters, and U.S. Pat. Nos. 4,433,496 and 5,964,547 disclose screw-driven wedges for securing adapters to buckets. While these devices eliminate the need for hammering, they each require a number of parts, thus, increasing the complexity and cost of the locks. The ingress of fines can also make removal difficult as the fines increase friction and interfere with the threaded connections. Moreover, with the use of standard threads, the fines can build up and become “cemented” around the threads to make turning of the bolt and release of the parts extremely difficult as can corrosion and damage to the threads.\nU.S. Pat. Nos. 6,986,216, 7,174,661 and 7,730,652 disclose locking arrangements for wear assemblies that rely upon a threaded wedge that engages a thread formation on the spool or wear member, and is rotated to drive the wedge into and out of the opening. These systems require minimal components, eliminate hammering, and alleviate the removal problems associated with prior systems. However, they lack the ability to provide substantial take up to ensure a tight fit with the lip or other supporting structure, or effective retightening after wear occurs.\nTypically, in a mining operation, a major earthmoving machine like a large cable shovel or dragline machine may have as many as three buckets dedicated to the machine. These buckets will include one bucket that is actively in use on the machine, one bucket that has been taken off the machine and is in the rebuild shop (e.g., to have various wear members removed and replaced with new wear members and to rebuild the lip for the tooth base and shroud fit areas), and one “ready line” bucket. The ready line bucket is a bucket that is new or has been through the re-build process and is ready to go back to work. The ready line bucket is needed because a bucket rebuild can take months to complete. It can be used on a scheduled maintenance cycle or, as can happen, when a major failure occurs with the bucket on the machine. Because the rebuild process takes so long, a mine cannot afford to not have a bucket available to put on a machine in case of emergency. The downtime and associated economic loss would be too great.\nWhile larger mining operations (e.g., operations involving multiple cable shovels and/or dragline machines) may not have three buckets dedicated to each machine, the operation will still typically have a sufficient number of ready line buckets available, if needed, to prevent excessive downtime (i.e., to avoid having a machine inoperable while waiting for a bucket rebuild job to be completed). The need for numerous ready line buckets represents a significant cost for the mining operation.\nBecause the lip rebuild tends to be the most time consuming part of the bucket rebuild process, reducing the number of rebuilds by lengthening the time between rebuilds would be a huge savings. Such a reduction in the number or frequency of rebuilds to the lip or other parts of the bucket would save the end user the money and time needed to perform these rebuilds as well as avoid the downtime associated with having the excavating bucket detached from the machine or unavailable for use in moving material. Reducing the number of lip rebuilds could constitute a huge savings in terms of less inventory of replacement buckets, fewer welders required to do these rebuilds, and a more forgiving system that is easier to operate and can be changed when it is more convenient for the operation.\nSince the bucket lip takes substantial abuse and is under considerable load during use, it needs to retain its strength and integrity to avoid failure. While welding on a lip rebuilds the leading edge of the lip to its original form, it also poses a risk to the lip if not done correctly. The lip must be preheated and welding procedures must be followed very carefully in order to avoid developing cracks. A cracked lip will necessitate the bucket being removed from the machine and repaired. However, if one does not need to weld repair the lip as often, then one possible failure mode is reduced or limited, thus minimizing the chances for a lip crack or failure.\nOne factor that may influence the need to repair or rebuild the lip on a bucket relates to whether the system for coupling the wear member to the lip is capable of securely engaging the parts together. The coupling system must be able to move the wear member a sufficient distance with respect to the lip to seat the wear member onto the lip. This amount of movement is referred to as “take up” (e.g., the coupling system must move the wear member a sufficient distance with respect to the lip to “take up” any gap or distance between the wear member and the lip). If a coupling system can only move a wear member a small distance with respect to the lip, the coupling system has a small take up capability, and in such systems, the mine operator may be forced to rebuild the lips more frequently (to assure that the coupling system will have sufficient take up to move the wear member and securely hold it against the lip). For coupling systems with a small amount of available take up, the lip rebuild also must be relatively precise to assure that the coupling system will be able to move the wear member and hold it onto the lip. Systems with wear members that are not tightly held to the supporting structure will tend to suffer more wear and tend to be more susceptible to wear member loss. While premature wearing of the lip may be of primary concern, premature wearing of other support structures, such as adapters, can also increase downtime and costs due to more frequent replacement.\nAccordingly, improvements in releasable coupling systems for securing wear members to the digging edge of a bucket would be welcome in the mining and construction industries. There remains a need for coupling systems that are easy and safe to install and remove, are reliable in use, enable substantial take up, allow longer time periods between bucket rebuilds, permit a wider range of dimensional variation in the manufacturing processes for the various parts, and lead to less machine downtime. Such improvements would result in reduced costs by decreasing the need for ready line buckets and the expense associated with rebuilding the digging edge of the buckets."}
{"text": "The present invention relates generally to electronic packaging. More particularly, the present invention relates to a marked structure such as a wafer or an array of packages.\nAs is well known to those of skill in the art, integrated circuits, i.e., electronic components, are fabricated in an array on a wafer. The wafer is then cut, sometimes called diced, to singulate the integrated circuits from one another.\nFIG. 1 is a cross-sectional view of a section of a wafer 10 being cut from a front-side surface 10F of wafer 10 in accordance with the prior art. Formed in wafer 10 were integrated circuits 12. Integrated circuits 12 were delineated by scribe lines 14, which included a first scribe line 14A and a second scribe line 14B, on front side surface 10F of wafer 10. For example, scribe lines 14 were formed by selective etching of a silicon oxide layer 18 on front-side surface 10F.\nTo illustrate, first scribe line 14A delineated a first integrated circuit 12A from a second integrated circuit 12B. Each scribe line 14 had a width WF.\nA back-side surface 10B of wafer 10 was attached to a tape 20. Wafer 10 was then sawed with a saw blade 22. Saw blade 22 was aligned with scribe lines 14 using an optical alignment system in a well-known manner. Saw blade 22 cut through wafer 10 along scribe lines 14. In this manner, integrated circuits 12 were singulated. Tape 20 supported wafer 10 during sawing and supported the singulated integrated circuits 12 after sawing was complete.\nGenerally, width WF of scribe lines 14 was sufficient to accommodate the width of the saw cut plus tolerance in the positioning of saw blade 22. Stated another way, width WF of scribe lines 14 was sufficiently large such that the saw cut made by saw blade 22 was always within a scribe line 14. For example, saw blade 22 is within scribe line 14B in FIG. 1.\nSince the optical alignment system used scribe lines 14 directly to align saw blade 22, saw blade 22 was aligned to scribe lines 14 to within tight tolerance. Accordingly, scribe lines 14 were relatively narrow and, more particularly, were only slightly wider than saw blade 22. To illustrate, width WF was within the range of 0.002 inches (0.051 mm) to 0.008 inches (0.203 mm).\nIn certain instances, it was important to protect the front-side surface of the wafer during sawing, e.g., from shards and particulates generated during sawing. To protect the front-side surface, the wafer was sawed from the back-side surface of the wafer as discussed below in reference to FIG. 2.\nFIG. 2 is a cross-sectional view of a section of a wafer 30 being cut from a back-side surface 30B of wafer 30 in accordance with the prior art. To protect a front-side surface 30F of wafer 30, front-side surface 30F was attached to a tape 32. Tape 32 supported wafer 30 during sawing.\nSaw blade 22 was aligned with scribe lines 14-1 on front-side surface 30F of wafer 30 using a two-step process. First, tape 32 was aligned with scribe lines 14-1. Front-side surface 30F was attached to tape 32. Tape 32 had area greater than the area of front-side surface 30F such that tape 32 had an exposed region, which extended beyond wafer 30. Tape 32 had alignment marks in the exposed region of tape 32. As an example, see alignment holes 30a and 30b of Roberts, Jr. et al., U.S. Pat. No. 5,362,681, which is herein incorporated by reference in its entirety. In the above manner, scribe lines 14-1 were aligned with the alignment marks of tape 32.\nSecond, saw blade 22 was aligned with the alignment marks of tape 32. Wafer 30 was then sawed with saw blade 22 from back-side surface 30B. However, since saw blade 22 was aligned indirectly to scribe lines 14-1 using alignment marks of tape 32, a large tolerance was associated with the alignment of saw blade 22 to scribe lines 14-1.\nTo accommodate this large tolerance, each of scribe lines 14-1 had a relatively large width WB. More particularly, referring now to FIGS. 1 and 2 together, width WB of scribe lines 14-1 of wafer 30, which was designed to be cut from back-side surface 30B, was significantly larger than width WF of scribe lines 14 of wafer 10, which was designed to be cut from front-side surface 10F. To illustrate, width WB was approximately 0.012 inches (0.305 mm) or more.\nDisadvantageously, forming scribe lines 14-1 with a relatively large width WB resulted in less integrated circuits 12 for any given size wafer 30 than the corresponding number of integrated circuits 12 formed in the same size wafer 10, i.e., there was a loss of yield of integrated circuits 12 from wafer 30. As a result, the cost of each integrated circuit 12 from wafer 30 was greater than the cost of each integrated circuit 12 from wafer 10. However, it is desirable to minimize the cost of each integrated circuit 12.\nIn accordance with the present invention, a method includes identifying and determining a position of a scribe grid on a front-side surface of a wafer with a camera. Based on this information, a computer aims a laser at a first location on a back-side surface of the wafer. The laser is fired to form a first alignment mark on the back-side surface of the wafer. Advantageously, the alignment mark is positioned with respect to the scribe grid to within tight tolerance, e.g., to within 0.001 inches (0.025 millimeters) or less.\nThe front-side surface of the wafer is attached to a tape to protect the front-side surface of the wafer during sawing. A saw blade is aligned with a scribe line of the scribe grid using the alignment mark on the backside surface of the wafer. The wafer is cut from the back-side surface along the scribe line with the saw blade.\nAdvantageously, the wafer is cut from the back-side surface thus protecting the front-side surface of the wafer and, more particularly, the integrated circuits. Of further importance, the saw blade is precisely aligned to the scribe line using the alignment mark such that the scribe line is not fabricated with the extra large width of scribe lines of conventional wafers designed to be cut from the back-side surface.\nRecall that in the prior art, in certain instances, it was important to cut the wafer from the back-side surface. However, to accommodate the large tolerance associated with back-side wafer cutting, the wafer designed to be cut from the back-side surface was formed with relatively wide scribe lines. Disadvantageously, forming the scribe lines with a relatively large width resulted in less integrated circuits for any given size wafer, i.e., a loss of yield. This resulted in a substantial increase in the cost of the integrated circuits.\nIn stark contrast, the wafer is cut from the back-side surface in accordance with the present invention without the associated loss of yield of the prior art. As a result, the integrated circuits of the wafer are protected during singulation yet are fabricated without the associated substantial increase in cost of the prior art.\nIn accordance with another embodiment of the present invention, an array of packages is marked. In accordance with this embodiment, a back-side surface of the array is scanned by a camera to identify and determine the position of fiducials on the back-side surface. Based on this information, a computer aims a laser at a first location on a front-side surface of the array. The laser is fired to form an alignment mark on the front-side surface of the array. Advantageously, the alignment mark is positioned with respect to the fiducials to within tight tolerance, e.g., to within 0.001 inches (0.025 millimeters) or less.\nThe back-side surface of the array is attached to a tape. A saw blade is aligned with the array using the alignment mark as a reference. The array is cut with the saw blade thus singulating the packages.\nA pick and place machine removes the packages from the tape. Advantageously, the packages are directly removed from the tape by a standard and relatively simple pick and place machine. Accordingly, removal of the packages from the tape is relatively simple and thus low cost. As a result, the packages are fabricated at a low cost.\nIn the prior art, an array of packages was singulated from the back-side surface. More particularly, the array was placed upside down on the tape such that a layer of encapsulant of the array was adhered to the tape and the fiducials extended upwards and were exposed. The array was singulated by cutting from the back-side surface using the fiducials as a reference.\nHowever, after singulation, the singulated packages had to be removed from the tape and inverted, e.g., using a pick and place machine with flip capability. The singulated packages had to be removed from the tape and inverted so that the singulated packages could be loaded into the grid carrier with the contacts (or other interconnection structure) facing downwards into the grid carrier. Back-end processing, e.g., automated attachment to the printed circuit mother board or automated testing, required that the singulated packages be loaded into the grid carrier in this manner. Disadvantageously, removing the singulated packages from the tape and inverting the packages required complex machinery, was labor intensive and, accordingly, increased the cost of the packages. In contrast, removal of the packages from the tape in accordance with the present invention is relatively simple and thus low cost.\nIn accordance with one particular embodiment, a structure includes a substrate, e.g., a wafer, having a first surface and a second surface. The structure further includes a reference feature, e.g., a scribe grid, on the first surface and at least one alignment mark on the second surface. The alignment mark has a positional relationship to the reference feature.\nIn accordance with another embodiment of the present invention, an array of packages includes a substrate having a first section, an electronic component such as an integrated circuit attached to the first section, metallizations on a first surface of the first section, and contacts on the metallizations. Bond pads of the integrated circuit are electrically connected to the contacts by bond wires. A layer of encapsulant covers the integrated circuit, the bond pads, the bond wires, the contacts, and the metallizations. A fiducial is on a first surface of the array and an alignment mark is on a second surface of the array. More particularly, the alignment mark is in the layer of encapsulant.\nIn one embodiment, a method includes identifying a reference feature on a first surface of a substrate such as a wafer or an array of packages. The method further includes marking a first location on a second surface of the substrate with a first alignment mark. The first alignment mark is used to determine a position of the reference feature.\nThese and other features and advantages of the present invention will be more readily apparent from the detailed description set forth below taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates generally to the field of illumination, and, more particularly, to a submersible color light. Although the present invention is subject to a wide range of applications, it is especially suited for use in a pool lighting system, and will be particularly described in that context.\nPool lights illuminate the water at night for the safety of swimmers and for aesthetic purposes. The illumination emanates from underwater lights affixed to the wall of the pool. As used herein, a pool is used generically to refer to a container for holding water or other liquids. Examples of such containers are recreational swimming pools, spas, and aquariums.\nTo enhance the aesthetics, some current underwater pool lights use a transparent color filter or shade affixed to the front of the lens of the pool light to filter the light emanating from the lens of the pool light and thus add color to the pool. The color filters come in a variety of colors but only one of these color filters can be affixed to the pool light at a given time. Thus, the color of the pool stays at that particular color that the color filter passes. In order to change the color of the pool, the color filter must be removed from the pool light and a different color filter installed across the lens of the pool light.\nAs a alternative to these fixed colored filters, a system has been devised whereby a rotating wheel having filters of several colors is provided, such as the system disclosed in U.S. Pat. No. 6,002,216 and incorporated herein by reference. In this arrangement, white light is provided from a single source to at least one fiber optic lens through an optical fiber. Colored light is emitted from each fiber optic lens by passing white light through the color filter wheel which is selectively rotated by a motor in the illuminator. The color of light emitted by multiple illuminators is synchronized by independent circuitry within each illuminator that responds to digital signals in the form of manually interrupted supply current.\nHowever, fiber optic underwater illumination systems have several limitations that lead to the need for the present invention. The first is that their performance is relative to the skill of the installer. Only a well-trained technician is capable of installing a fiber optic system that can adequately illuminate a swimming pool. The availability of qualified training is limited thus the availability of trained installers is limited. Rushed fiber termination or fiber termination performed by an untrained installer can result in more than a 30% decrease in fiber optic system performance and can ultimately result in a costly failure of the total fiber optic system.\nThe second disadvantage of underwater fiber optic illumination is the limited amount of light delivered to the pool. This results from the light attenuation over distance that is inherent in the fibers' composition and the inefficiencies of focusing available light into the optical fiber at the light source.\nA further drawback of fiber optic underwater illumination is in the possibility of retrofitting the millions of existing pools having traditional submersible incandescent lighting fixtures. The feasibility of installing adequately sized fiber optic cable in the existing conduits is limited. These conduits are commonly ½ inch in diameter and are rarely over one inch in diameter. The minimum conduit diameter to carry a single fiber optic cable capable of delivering minimally acceptable light to a pool is one inch and the recommended size is 1½ inches.\nAn additional limitation of fiber optic systems is the additional cost of the materials and professional installation.\nThe alternative to colored fiber optic systems, providing colored lenses to submersible incandescent lighting fixtures, can be troublesome as well. These fixtures can be supplied with a colored glass lens to deliver that specific color to the pool. These colored glass lenses are typically limited to how richly they can color the light because the darker (or richer) the lens color, the more light in the form of heat that is trapped in the lens and the fixture. As the lens becomes too hot by absorbing too much light it can break due to thermal expansion or due to the differences in thermal expansion on the hot interior surface of the glass and the cool exterior surface that is in contact with the water. Further, as a result less light is emitted and it may be insufficient to illuminate the pool.\nAs an alternative to glass lenses, snap on or twist lock plastic colored lenses can be installed over an existing clear glass lens for a considerably simpler method to changing the color of the pool lighting. This method still requires physically lying or kneeling on the edge of the pool an reaching below the water to remove the existing plastic lens and then reaching again into the water to install the next colored plastic lens. Economical transparent colored plastics are also inefficient light transmitters reducing the amount of colored light reaching the pool.\nA need therefore exists for pool lights that can easily replace existing self-contained, incandescent lighting fixtures, but having synchronized color wheels without the additional cost of installing fiber optic cables and other drawbacks associated with fiber optic underwater illumination systems. Further, a need exists for colored lenses to be used with incandescent fixtures that do not trap excessive amounts of light and heat."}
{"text": "1. Field of the Art\nEmbodiments of the present invention relate to electronic devices for viewing stereographic displays and, in particular, to actively shuttered three-dimensional (“3-D” or “3D”) glasses synchronized in time with 3D enabled televisions.\n2. Description of the Related Art\nTelevisions (TVs), monitors, and other displays have evolved over the years from black and white images to color, analog signals to digital high definition television (HDTV) formats, and cathode ray tube (CRT) screens to plasma and liquid crystal display (LCD) technologies. Three-dimensional ready televisions may be the next major upgrade that the general public endorses by buying in large quantities.\nThree-dimensional ready televisions commonly incorporate stereoscopic displays. Stereoscopic displays present a separate image to a viewer's left eye and a separate image to the viewer's right eye. Technologies to present separate images to a viewer's eyes include anaglyphs, which typically use red and blue lens glasses, polarized lens glasses, and active shutter lens glasses for actively blocking the viewer's eyes in rapid, period sequence. All such lenses for 3D glasses are typically non-corrective lenses in that they do not redirect the direction of light (e.g., by refraction) as it passes through the lens.\nFor many 3D-ready televisions coming to market, their manufactures have developed vendor-unique protocols for viewing based on active shutter techniques. With active shutter techniques, a 3D television rapidly (e.g., 30 frames per second or greater) and alternatingly shows separate images for a viewer's left and right eyes. A viewer typically wears glasses that have liquid crystal (LC) active shutters rapidly blocking the left and right eye views alternatively so that each eye sees the corresponding left and right images shown at the corresponding time on the display. This “active shutter” process preferably is periodically and/or continuously synchronized with synchronization signals transmitted from the TV.\nCurrently, the synchronization signal is broadcasted from the TV wirelessly via a specialized transmitter in the television either built-in inside the TV front panel or external as peripheral, such as an infrared (IR) transmitter. The glasses typically comprise a specialized receiver to receive and analyze the wirelessly transmitted synchronization signal. This synchronization signal aligns the 3D glasses' lens timing to that of the television and corrects for wander in the 3D glasses' clock circuitry."}
{"text": "Bacterial cells, such as E. coli, are commonly used for producing recombinant proteins. There are many advantages to using bacterial cells, such as E. coli, for producing recombinant proteins particularly due to the versatile nature of bacterial cells as host cells allowing the gene insertion via plasmids. E. coli have been used to produce many recombinant proteins including human insulin.\nDespite the many advantages to using bacterial cells to produce recombinant proteins, there are still significant limitations including the difficulty of producing protease sensitive proteins. Proteases play an important role in turning over old, damaged or miss-folded proteins in the E. coli periplasm and cytoplasm. Bacterial proteases act to degrade the recombinant protein of interest, thereby often significantly reducing the yield of active protein.\nA number of bacterial proteases have been identified. In E. coli proteases including Protease III (ptr), DegP, OmpT, Tsp, prlC, ptrA, ptrB, pepA-T, tsh, espc, eatA, clpP and lon have been identified.\nTsp (also known as Prc) is a 60 kDa periplasmic protease. The first known substrate of Tsp was Penicillin-binding protein-3 (PBP3) (Determination of the cleavage site involved in C-terminal processing of penicillin-binding protein 3 of Escherichia coli; Nagasawa H, Sakagami Y, Suzuki A, Suzuki H, Hara H, Hirota Y. J. Bacteriol. 1989 November; 171(11):5890-3 and Cloning, mapping and characterization of the Escherichia coli Tsp gene which is involved in C-terminal processing of penicillin-binding protein 3; Hara H, Yamamoto Y, Higashitani A, Suzuki H, Nishimura Y. J. Bacteriol. 1991 August; 173 (15):4799-813) but it was later discovered that the Tsp was also able to cleave phage tail proteins and, therefore, it was renamed as Tail Specific Protease (Tsp) (Silber et al., Proc. Natl. Acad. Sci. USA, 89: 295-299 (1992)). Silber et al. (Deletion of the prc(tsp) gene provides evidence for additional tail-specific proteolytic activity in Escherichia coli-12; Silber, K. R., Sauer, R. T.; Mol Gen Genet 1994 242:237-240) describes a prc deletion strain (KS1000) wherein the mutation was created by replacing a segment of the prc gene with a fragment comprising a Kanr marker.\nThe reduction of Tsp (prc) activity is desirable to reduce the proteolysis of proteins of interest. However, it was found that cells lacking protease prc show thermosensitive growth at low osmolarity. Hara et al isolated thermoresistant revertants containing extragenic suppressor (spr) mutations (Hara et al., Microbial Drug Resistance, 2: 63-72 (1996)). Spr is an 18 kDa membrane bound periplasmic protease and the substrates of spr are Tsp and peptidoglycans in the outer membrane involved in cell wall hydrolysis during cell division. The spr gene is designated as UniProtKB/Swiss-Prot P0AFV4 (SPR_ECOLI).\nImproved protease deficient strains comprising mutant spr gene have been described. Chen et al describes the construction of E. coli strains carrying different combinations of mutations in prc (Tsp) and another protease, DegP, created by amplifying the upstream and downstream regions of the gene and ligating these together on a vector comprising selection markers and a sprW174R mutation (High-level accumulation of a recombinant antibody fragment in the periplasm of Escherichia coli requires a triple-mutant (ΔDegP Δprc sprW174R) host strain (Chen C, Snedecor B, Nishihara J C, Joly J C, McFarland N, Andersen D C, Battersby J E, Champion K M. Biotechnol Bioeng. 2004 Mar. 5; 85(5):463-74). The combination of the ΔDegP, Δprc and sprW174R mutations were found to provide the highest levels of antibody light chain, antibody heavy chain and F(ab′)2-LZ. EP1341899 discloses an E. coli strain that is deficient in chromosomal DegP and prc encoding proteases DegP and Prc, respectively, and harbors a mutant spr gene that encodes a protein that suppresses growth phenotypes exhibited by strains harboring prc mutants.\nThe present invention provides new bacterial strains carrying alternative spr mutants which provide advantageous means for producing recombinant proteins."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor storage device and a semiconductor storage device driving method.\n2. Related Art\nNonvolatile semiconductor storage devices such as NAND flash memories are frequently used in portable devices such as digital cameras and cellular phones. Such a semiconductor storage device is required to be low in power dissipation, low in current dissipation, and fast in operation.\nSuch a semiconductor storage device includes a memory cell array formed on a well in a semiconductor substrate. The memory cell array is divided into a large number of blocks, and each block includes a plurality of memory cells. One well is formed so as to extend over a plurality of blocks. Typically, one well is provided so as to extend over the entire memory cell, When erasing data in memory cells included in one block, therefore, it is necessary to charge a large-capacity well provided over a plurality of blocks. Hereafter, a well potential needed to erase data in memory cells is referred to as erasing potential.\nThe well is boosted by using booster circuits of a plurality of kinds sometimes, For example, the well is boosted by booster circuits of two kinds: a booster circuit BCpgm and a booster circuit BCpass. The booster circuit BCpgm supplies a write voltage to memory cells into which data is to be written, i.e., selected cells The booster circuit BCpass supplies a write inhibition voltage to memory cells into which data is not to be written, i.e., inhibited memory cells. The booster circuit BCpgm has more booster stages than the booster circuit BCpass. As a result, the booster circuit BCpgm can conduct boosting to a higher potential than the booster circuit BCpass. On the other hand, the booster circuit BCpass is less in power dissipation than the booster circuit BCpgm.\nIn order to boost the well to an erasing potential at the time of data erasing, either a method using only the booster circuit BCpgm (method 1) or a method using both the booster circuits BCpgm and BCpass (method 2) is conventionally used.\nIn the method it takes a considerably long time to boost the large-capacity well to the erasing potential. Therefore, it is difficult for the method 1 to meet the demand of fast operation of semiconductor storage devices which will become more intensive hereafter. In the method 2 boosting is conducted by using both booster circuits BCpgm and BCpass, and then boosting is conducted by using only the booster circuit BCpgm. In the method 2 therefore, the well can be boosted in a short time. In the method 2, however, both the booster circuits BCpgm and BCpass are used, resulting in a problem that the instantaneous peak of the current dissipation becomes too large.\nIn both the method 1 and the method 2, the booster circuit BCpgm is mainly used to boost the well. Since the booster circuit BCpgm has more booster stages than BCpass, the booster circuit BCpgm is poor in boosting efficiency. Therefore, these methods also have a problem that considerably much power is dissipated to boost the well\nTherefore, a semiconductor storage device and a semiconductor storage device driving method that are low in power dissipation, low in current dissipation, and fast in operation speed are provided."}
{"text": "1. Field of the Invention\nThe invention is in the field of biotechnology. Specifically, the invention is directed to compositions of mutated binding proteins containing reporter groups, analyte biosensor devices derived therefrom, and their use as analyte biosensors both in vitro and in vivo.\n2. Description of Relevant Art\nMonitoring glucose concentrations to facilitate adequate metabolic control in diabetics is a desirable goal and would enhance the lives of many individuals. Currently, most diabetics use the “finger stick” method to monitor their blood glucose levels and patient compliance is problematic due to pain caused by frequent (several times per day) sticks. As a consequence, there have been efforts to develop non-invasive or minimally invasive in vivo and more efficient in vitro methods for frequent and/or continuous monitoring of blood glucose or other glucose-containing biological fluids. Some of the most promising of these methods involve the use of a biosensor. Biosensors are devices capable of providing specific quantitative or semi-quantitative analytical information using a biological recognition element that is combined with a transducing (detecting) element.\nThe biological recognition element of a biosensor determines the selectivity, so that only the compound to be measured leads to a signal. The selection may be based on biochemical recognition of the ligand where the chemical structure of the ligand (e.g. glucose) is unchanged, or biocatalysis in which the element catalyzes a biochemical reaction of the analyte.\nThe transducer translates the recognition of the biological recognition element into a semi-quantitative or quantitative signal. Possible transducer technologies are optical, electrochemical, acoustical/mechanical or colorimetrical. The optical properties that have been exploited include absorbance, fluorescence/phosphorescence, bio/chemiluminescence, reflectance, light scattering and refractive index. Conventional reporter groups such as fluorescent compounds may be used, or alternatively, there is the opportunity for direct optical detection, without the need for a label.\nBiosensors specifically designed for glucose detection that use biological elements for signal transduction typically use electrochemical or colorimetric detection of glucose oxidase activity. This method is associated with difficulties including the influence of oxygen levels, inhibitors in the blood and problems with electrodes. In addition, detection results in consumption of the analyte that can cause difficulties when measuring low glucose concentrations.\nA rapidly advancing area of biosensor development is the use of fluorescently labeled periplasmic binding proteins (PBP's). As reported by Cass (Anal. Chem. 1994, 66, 3840-3847), a labeled maltose binding protein (MBP) was effectively demonstrated as a useable maltose sensor. In this work MBP, which has no native cysteine residues, was mutated to provide a protein with a single cysteine residue at a position at 337 (S337C). This mutation position is within the binding cleft where maltose binding occurred and therefore experiences a large environmental change upon maltose binding. Numerous fluorophores were studied, some either blocked ligand binding or interfered with the conformational change of the protein. Of those studied IANBD resulted in a substantial increase in fluorescence (160%) intensity upon maltose binding. This result is consistent with the location of the fluorophore changing from a hydrophilic or solvent exposed environment to a more hydrophobic environment as would have been theoretically predicted for the closing of the hinge upon maltose binding. However, this mutant protein and the associated reporter group do not bind diagnostically important sugars in mammalian bodily fluids. Cass also disclosed (Analytical Chemistry 1998, 70(23), 5111-5113) association of this protein onto TiO2 surfaces, however, the surface-bound protein suffered from reduced activity with time and required constant hydration.\nHellinga, et al. (U.S. Pat. No. 6,277,627) reports the engineering of a glucose biosensor by introducing a fluorescent transducer into a Galactose/Glucose Binding Protein (GGBP) mutated to contain a cysteine residue, taking advantage of the large conformation changes that occur upon glucose binding. Hellinga, et. al. (U.S. Pat. No. 6,277,627) disclose that the transmission of conformational changes in mutated GGBPs can be exploited to construct integrated signal transduction functions that convert a glucose binding event into a change in fluorescence via an allosteric coupling mechanism. The fluorescent transduction functions are reported to interfere minimally with the intrinsic binding properties of the sugar binding pocket in GGBP.\nIn order to accurately determine glucose concentration in biological solutions such as blood, interstitial fluids, ocular solutions or perspiration, etc., it may be desirable to adjust the binding constant of the sensing molecule of a biosensor so as to match the physiological and/or pathological operating range of the biological solution of interest. Without the appropriate binding constant, a signal may be out of range for a particular physiological and/or pathological concentration. Additionally, biosensors may be configured using more than one protein, each with a different binding constant, to provide accurate measurements over a wide range of glucose concentrations as disclosed by Lakowicz (U.S. Pat. No. 6,197,534).\nDespite the usefulness of mutated GGBPs, few of these proteins have been designed and examined, either with or without reporter groups. Specific mutations of sites and/or attachment of certain reporter groups may act to modify a binding constant in an unpredictable way. Additionally, a biosensor containing reporter groups may have a desirable binding constant, but not result in an easily detectable signal change upon analyte binding. Some of the overriding factors that determine sensitivity of a particular reporter probe attached to a particular protein for the detection of a specific analyte is the nature of the specific interactions between the selected probe and amino acid residues of the protein. It is not currently possible to predict these interactions within proteins using existing computational methods, nor is it possible to employ rational design methodology to optimize the choice of reporter probes. It is currently not possible to predict the effect on either the binding constant or the selectivity based on the position of any reporter group, or amino acid substitution in the protein (or vice-versa).\nTo develop reagentless, self-contained, and/or implantable and/or reusable biosensors using proteins the transduction element must be in communication with a detection device to interrogate the signal to and from the transduction element. By “interrogate” is meant the process of sending light to, and measuring emitted light from, a sample. Typical methods include placing proteins within or onto the surface of optical fibers or planner waveguides using immobilization strategies. Such immobilization strategies include entrapment of the protein within semi-permeable membranes, organic polymer matrixes, or inorganic polymer matrixes. The immobilization strategy may ultimately determine the performance of the working biosensor. The prior art details numerous problems associated with the immobilization of biological molecules. For example, many proteins undergo irreversible conformational changes, denaturation, and loss of biochemical activity. Immobilized proteins can exist in a large number of possible orientations on any particular surface, for example, with some proteins oriented such that their active sites are exposed whereas others may be oriented such that there active sites are not exposed, and thus not able to undergo selective binding reactions with the analyte. Immobilized proteins are also subject to time-dependent denaturation, denaturation during immobilization, and leaching of the entrapped protein subsequent to immobilization. Therefore, problems result including an inability to maintain calibration of the sensing device and signal drift. In general, binding proteins require orientational control to enable their use, thus physical absorption and random or bulk covalent surface attachment or immobilization strategies as taught in the literature generally are not successful.\nTherefore, there is a need in the art to design additional useful mutated proteins and mutated GGBP proteins generating detectable signal changes upon analyte binding for use as biosensors, and additionally there is a need in the art to design additional useful mutated binding protein and mutated GGBPs containing reporter groups generating detectable and reversible signal changes upon analyte or glucose binding for use as biosensors."}
{"text": "NAND flash memory, as well as other types of non-volatile memories (“NVMs”), are commonly used for mass storage. For example, consumer electronics such as portable media players often include flash memory to store music, videos, and other media.\nTypically, each page of a NVM can be used to store a single logical block address (“LBA”). However, when the system needs to access smaller pieces of data (e.g., a portion of a page), an entire page has to be accessed. For example, during a read operation, the system may need to read an entire page and discard the unwanted data."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nTrucks such as pickups and trucks having a cargo box can be a popular type of vehicle because the cargo box enables objects of various sizes to be placed therein and transported. The cargo box pickup generally has a floor surface or bed that is bounded by a forward wall, opposing sidewalls and a tailgate providing a rearward wall.\nCovers can be used to enclose the cargo box to protect the contents against dirt, debris, and other environmental contaminants, and to improve the aesthetic quality thereof. In some cases, the covers can be formed of generally planar sections or panels hingedly coupled together so they can be folded together in a stacked arrangement to allow access to the cargo box. The stacked panels typically still leave a portion of the cargo box covered.\nIn some cases the stacked panels can be rotated so the tonneau cover is angled and extends upwardly above the truck box. Prop-rods can be used to retain the stacked panels in their upwardly angled position, but can require multiple steps, including coupling components together, and/or can limit locking options or angular orientations."}
{"text": "Motorized marine vehicles such as boats generally utilize an internal fuel tank. The fuel tank of the marine vehicle typically communicates with a neck or fuel conduit having an opening which opens to the exterior of the marine vehicle and through which a fuel nozzle is inserted to fill the fuel tank with fuel. A vent stack or overflow conduit may extend from the fuel tank and discharge to the side of the marine vehicle to vent fumes from the fuel tank. In the event that the fuel tank is overfilled with fuel, a visual indication of such overfilling may result from flow of fuel initially from the fuel tank through the overflow conduit and ultimately, discharge of the fuel from the overflow conduit onto the water body. However, overfilling of the fuel tank and discharge of the fuel into the water body may contribute to pollution of the water body. Furthermore, the discharged fuel may become a potential fire hazard since the fuel will float on the water and remain combustible. Additionally, air may become blocked in the overflow conduit and fuel may spill from the conduit even though the fuel tank is not full. Thus, knowledge of when to stop the flow of fuel into the fuel tank on the part of the person filling the tank to prevent tank overfilling is desirable.\nTherefore, a need exists for a fuel overflow alarm system which indicates that overfilling of a fuel tank in a marine vehicle is imminent during refueling of the marine vehicle to prevent fuel from being discharged into a water body on which the marine vehicle is docked."}
{"text": "The present invention relates generally to intraoral devices and, more particularly, to an intraoral discluder for use in relieving tension headaches, common migraine headaches, and temporomandibular disorders.\nTension and muscle contraction headaches affect many people every day. The headaches are often recurring and, without effective treatment, can become very painful, restricting an individual's ability to think clearly and function effectively. The discomfort associated with tension and muscle contraction headaches is usually due to pain from strained and fatigued muscles of the head. The majority of the muscles of the human head are not sufficiently strong to elicit the type of pain and discomfort associated with tension and muscle contraction headaches. That is not the case with the temporalis muscle, however, which is located on the side of the skull and extends from just behind the eye to just behind the ear, and which is an extremely powerful muscle that functions to close or elevate the jaw.\nUnder normal circumstances, the temporalis muscle should not exert a large static force by contracting isometrically, except possibly during normal chewing. Inappropriate isometric contraction of the temporalis muscle is commonly known as “clenching” and is clinically known as myofascial dysfunction. The intensity of the myofascial dysfunction varies according to the mandible's anterior/posterior position, with the intensity increasing as the mandible's position moves posteriorly. Unfortunately, myofascial dysfunction is particularly difficult to detect or diagnose, because the act of clenching is a relatively motionless act that is commonly done while a person is concentrating on another topic, or while sleeping.\nAs the muscular contraction condition of “clenching” continues, the muscle becomes fatigued and susceptible to spasm and cramping. The pain from spasming and cramping temporalis fibers is severe and is usually diagnosed as a common migraine. Headache sufferers who seek the assistance of a physician typically are treated with muscle relaxants, analgesics, and/or physical therapy for the muscle fatigue. However, medications and physical therapy require continual treatment, and they treat only the symptoms of the underlying problem, not the source of the problem itself.\nHeadache sufferers who seek the assistance of a dentist typically are diagnosed as having a temporomandibular disorder and are treated with an intraoral “jaw-positioning” appliance. Unfortunately, the intraoral appliances provided by dentists frequently are not entirely effective, because they only approximate the relative positions of the upper and lower teeth with respect to each other, allowing clenching to continue with minimal mandibular movement. Further, these intraoral appliances ordinarily cannot be used by patients who have malocclusions, protrusions or retrusions of the mandible, or other irregular teeth or mandibular orientations. Typically, the intraoral appliance must also be fabricated by a dentist at a prohibitive cost to a majority of individuals who suffer from tension headaches and common migraine headaches. Lastly, most intraoral jaw-positioning appliances and other types of semi-custom intraoral discluders can be used only on the upper teeth. However, in some circumstances, use of the appliance on the upper teeth is impossible due to malocclusions and irregular orientation of the teeth.\nOne intraoral appliance that avoids the drawbacks mentioned immediately above is disclosed in U.S. Pat. No. 5,795,150 to Boyd, Sr. That appliance includes a trough sized and configured to be releasably retained by a wearer's maxillary incisors and further includes a dome projecting posteriorly from the trough and defining a surface to be contacted by at least one opposing incisor. When the appliance is properly positioned in the wearer's mouth, the temporalis muscles are rendered ineffective, thus relieving tension headaches, common migraine headaches, and temporomandibular disorders. However, it is believed that this appliance can sometimes still allow limited clenching of the temporalis muscle, particularly when the mandible is located in its furthest posterior position.\nIt should be apparent from the foregoing discussion, that there remains a need for an even more effective intraoral discluder configured to be placed on either the upper teeth or the lower teeth, to prevent contact of the upper and lower teeth in all mandibular movements and to further inhibit undesired isometric contraction of the temporalis muscle. The present invention satisfies this need."}
{"text": "1. Field of the Invention\nThe present invention relates to a phase shifter, in more particular, to a transmission line phase shifter.\n2. Related Art\nAs a phase shifter used for beam control or phase modulation of a phased array antenna, a transmission line phase shifter has been conventionally used. For example, Japanese Patent Laid-Open No. 5-14004 discloses a phase adjustment circuit, comprising a first substrate, a U-shaped pattern formed on the first substrate, a second substrate, and first and second patterns formed on the second substrate, each of the first and second patterns having a part provided in parallel with each other, in which parallel parts in the U-shaped pattern are overlapped with the parallel patterns of the first and second patterns to contact to each other, and the first and second substrates are configured to be continuously movable.\nAccording to the phase adjustment circuit disclosed by Japanese Patent Laid-Open No. 5-14004, a length of the U-shaped pattern is determined as a length of an integral multiplication of a ½ wavelength of a signal to be transmitted, so that the first and second substrates are continuously movable in a state where the respective parallel parts in the U-shaped pattern are overlapped with and in contact with the parallel patterns of the first and second patterns. According to this structure, in the phase adjustment circuit of Japanese Patent Laid-Open No. 5-14004, it is possible to continuously change a transmission path length of the signal, so as to continuously change a signal phase while confirming circuit characteristics.\nIn addition, Japanese Patent Laid-Open No. 2001-237605 discloses a phase shifter, comprising a first dielectric substrate, a plurality of input side microstrip lines and a plurality of output side microstrip lines provided on the first dielectric substrate, a second dielectric substrate which is movable with respect to the first dielectric substrate, a plurality of coupling microstrip lines provided on the second dielectric substrate, and an insulator provided between the first dielectric substrate and the second dielectric substrate, in which a plurality of the input side microstrip lines and a plurality of the output side microstrip lines are overlapped with a plurality of the coupling microstrip lines respectively to be facing to each other.\nAccording to the phase shifter disclosed by Japanese Patent Laid-Open No. 2001-237605, it is possible to change lengths of overlapped portions, in which a plurality of the input side microstrip lines and a plurality of the output side microstrip lines are overlapped with a plurality of the coupling microstrip lines via the insulator, simultaneously with a constant proportion. According to this structure, it is possible to change phases of the signals transmitted through the input side microstrip lines the respective coupling microstrip lines at the same time. For example, the phase shifter disclosed by Japanese Patent Laid-Open No. 2001-237605 may be installed in an array antenna used for an antenna for a portable telephone base station as an apparatus for changing a directivity orientation.\nHowever, in the phase adjustment circuit disclosed by Japanese Patent Laid-Open No. 5-14004, a total length of the U-shaped pattern is previously set (fixed) as the length of the integral multiplication of the ½ wavelength of the signal to be transmitted.\nIn the phase shifter disclosed by Japanese Patent Laid-Open No. 2001-237605, a total length of a plurality of the coupling microstrip lines is previously set (fixed) as the length of the integral multiplication of the ½ wavelength of the wavelength of the signal to be transmitted respectively.\nTherefore, in both of the phase adjustment circuit disclosed by Japanese Patent Laid-Open No. 5-14004 and the phase shifter disclosed by Japanese Patent Laid-Open No. 2001-237605, it is difficult to improve transmission characteristics and return loss characteristics when a signal at a frequency other than a frequency of the signal that is supposed to be used when designing the phase shifter."}
{"text": "This invention is directed to α-hydroxyamide compounds. In particular, this invention is directed to α-hydroxyamide compounds that are bradykinin antagonists or inverse agonists.\nBradykinin (“BK”) is a kinin which plays an important role in the pathophysiological processes accompanying acute and chronic pain and inflammation. Bradykinin (BK), like other kinins, is an autacoid peptide produced by the catalytic action of kallikrein enzymes on plasma and tissue precursors termed kininogens. The biological actions of BK are mediated by at least two major G-protein-coupled BK receptors termed B1 and B2. It is generally believed that B2 receptors, but not B1 receptors, are expressed in normal tissues and that inflammation, tissue damage or bacterial infection can rapidly induce B1 receptor expression. This makes the B1 receptor a particularly attractive drug target. The putative role of kinins, and specifically BK, in the management of pain and inflammation has provided the impetus for developing potent and selective BK antagonists. In recent years, this effort has been heightened with the expectation that useful therapeutic agents with analgesic and anti-inflammatory properties would provide relief from maladies mediated through a BK receptor pathway (see e.g., M. G. Bock and J. Longmore, Current Opinion in Chem. Biol., 4:401-406 (2000)). Accordingly, there is a need for novel compounds that are effective in blocking or reversing activation of bradykinin receptors. Such compounds would be useful in the management of pain and inflammation, as well as in the treatment or prevention of diseases and disorders mediated by bradykinin; further, such compounds are also useful as research tools (in vivo and in vitro)."}
{"text": "1. Field of the Invention\nThe present disclosure relates to a device for determining dimensions of a workpiece, and particularly to a device for determining dimensions of a workpiece without having to contact the workpiece.\n2. Description of Related Art\nDimensions of many workpieces require determining. At present, simplified dimension determining tools (for example, vernier calipers and micrometers), and complicated determining apparatuses (for example, three-dimensional determining apparatuses or image detecting apparatuses), are all widely used to measure dimensions of workpieces.\nIn the dimension determining process, most dimension determining tools or apparatuses require contact with or are pressed towards workpieces. Therefore, surfaces or shapes of the workpieces may be damaged.\nTherefore, there is room for improvement within the art."}
{"text": "1. Field of the Invention\nThis invention relates to a reinforcing spacer device for use in the horizontal and vertical mortar of masonry walls. More particularly, this invention relates to a device combining reinforcement and spacing between adjacent masonry units.\n2. Description of the Prior Art\nMasonry, the building of structures from individual units laid in and bound together by mortar, is commonly used for the construction of buildings. Such widespread use is the result of the high durability, compressive strength, thermal mass and heat resistance of the masonry building materials. Because masonry construction requires extensive manual labor and individual building materials such as brick or block, the quality of the masonry construction is directly dependent on the type of materials used and the workmanship of the mason.\nIn the past, the construction of structurally sound, high quality wall units constructed of brick, block, stone or similar masonry building units has depended upon the workmanship of masons. The pattern used in setting the building materials strongly effects the durability of the overall construction. Through extensive training, masons develop techniques for evenly applying mortar to bed joints and vertical joints between units so that the units are evenly spaced in each course and the wall face remains vertical.\nIn recent years attention has been paid to wall reinforcement especially for areas that are routinely subjected to seismic forces. Here both wire reinforcements and ladder and truss reinforcements have been in widespread use. Such reinforcements are embedded in the horizontal mortar joints to reinforce, bond and control shrinkage cracking. Additionally, the reinforcements provide higher resistance to lateral loads, such as wind, by increasing tensile strength.\nThe inventors' patents and their assignee's product line include masonry accessories, namely, ladder and truss reinforcements, wall anchors, veneer ties, masonry flashing and related items for cavity walls. These products, which are sold under the trademarks of Lox All, DW-10X, X-seal and FlexFlash, are manufactured by Hohmann & Barnard, Inc., Hauppauge, N.Y. 11788 (“H&B”). The products have become widely accepted in the construction industry and the inventors have gained particular insight into the technological needs of this marketplace.\nIn general, the difficulties with masonry construction lie in two distinct areas. The first is the weakness of the horizontal mortar or bed joints that bond the masonry units together. This weakness is the result of the low tensile strength of the mortar joints and generally requires mortar joint reinforcement for structural stability. The second difficulty is constructing a wall with consistent and uniform mortar joints that keeps the structure vertical and maintains aesthetics. An uneven mortar joint thickness detracts from the overall appearance of the wall and can effect the overall stability and durability of the masonry construction.\nThe first difficulty is addressed by well-known devices, such as ladder and truss reinforcements for augmenting the tensile strength of the horizontal mortar joints. These devices greatly reduce cracking that can arise from thermal stresses, and increase lateral flexural strength, elasticity and performance of masonry walls under various stresses.\nExemplary of the above, in a patent to Stephen Priest, Jr., U.S. Pat. No. 903,000 issued Nov. 3, 1908, entitled “Wall Tie,” a reinforcing ladder device is taught which is constructed of twisted wires with one side of the ladder device embedded in the outer wythe and the other, in the inner wythe. Similarly, H. Spaight, U.S. Pat. No. 2,300,181 issued Oct. 27, 1942, entitled “Means for Constructing Buildings,” teaches a truss-shaped reinforcement device for embedment in either one wythe or in cavity walls in both wythes. More recently, W. Smith in U.S. Pat. No. 3,183,628 issued May 18, 1965, entitled “Masonry Wall Reinforcing Means,” describes an improvement of the Spaight invention by teaching truss and ladder reinforcements having grooves or bosses on the parallel side wires to increase the mortar bonding therewith. Placing one of the aforementioned devices in the horizontal mortar joints enhances the tensile strength of the horizontal joints.\nSeveral improvements to masonry wall reinforcement have been made by H&B. In 1976, Hala and Schwalberg of H&B, received U.S. Pat. No. 3,964,226 for an adjustable wall-tie reinforcing system which joined reinforcements in inner and outer wythes with an attached eye and pintle structure. During the period when the Uniform Building Code developed joint reinforcement specifications, Hohmann et al., received U.S. Pat. No. 5,454,200 issued Oct. 3, 1995 and U.S. Pat. No. 6,279,283 issued Aug. 28, 2001. These patents provide veneer anchoring systems for masonry walls which include reinforcement for cavity walls. These devices have received widespread usage in the industry. However, none of these devices were designed to aid in the application of uniformly thick mortar joints.\nIn the past, builders inserted devices between the block, stone or similar masonry building material to align and space the mortar joints. Representative patents include: U.S. Pat. No. 2,172,816—A. Douglas and V. Lefebure—Issued Sep. 12, 1939, entitled “Construction of Walls, Partitions, and the Like,” which describes a method of aligning specially-slotted or grooved building blocks using T-shaped dowel plates to tie three adjacent blocks; U.S. Pat. No. 4,334,397—G. Hitz—Issued Jun. 15, 1982, entitled “Masonry Structure and Apparatus and Process for Spacing Block in the Structure,” which describes a plastic-pronged spacer that separates the blocks; and U.S. Pat. No. 5,351,457—W. Colen—Issued Oct. 4, 1994, entitled “Wall Construction and Spacer for Use Therewith,” which describes wall construction spacers for use in various blocks that tie two adjacent blocks together.\nEach of the aforementioned devices addresses either reinforcement or spacing. Additionally, more recent technical advances have resulted in several combination reinforcing and spacing devices, described below. However, these devices do not provide a combination unit for both the reinforcing of the horizontal mortar joints and the spacing of the horizontal and vertical mortar joints.\nIn preparing for this application the below-mentioned patents have become known to the inventors hereof.\nPatentInventorIssue Date6,629,393PignataroOct. 7, 20036,553,737BergApr. 29, 20036,279,283Hohmann et al.Aug. 28, 20015,454,200Hohmann et al.Oct. 3, 19955,351,457ColenOct. 4, 19945,259,161CarterNov. 9, 19914,689,931HodgesSep. 1, 19874,334,397HitzJun. 15, 19823,964,226Hala et al.Jun. 22, 19763,183,628SmithMay 18, 19652,300,181SpaightOct. 27, 19422,172,816Douglas et al.Sep. 3, 1939903,000PriestNov. 3, 1908\nOf these patents, those not previously discussed, are discussed hereinbelow.\nInventions to properly space masonry building materials are described in patents such as P. Hodges in U.S. Pat. No. 4,689,931 issued Sep. 1, 1987, entitled “Masonry Construction Device,” and T. Berg in U.S. Pat. No. 6,553,737 issued Apr. 29, 2003, entitled “Method and apparatus to achieve consistent spacing between layers of modular construction material.”\nHodges describes a squared “U” shaped device with a straight bridging member connected between parallel sawtooth shaped members. Upon connecting adjacent ones of Hodges devices, the bridging members thereof form supports for the next course of masonry materials and thereby spaces consecutive courses. The “U” shaped devices are short and manually chained together to cover horizontally adjacent bricks. Berg teaches a device with spacer studs connected to the parallel and transverse rails of a ladder-type reinforcement. The device is placed on the top surface of a brick or block to space the horizontal mortar joint. Along a similar path, F. Carter in U.S. Pat. No. 5,259,161 issued Nov. 9, 1991, entitled “Vertical and horizontal reinforcement and spacing guide for panels constructed of blocks,” teaches, specifically for glass block installation, a grid formed by detailed thick elongate reinforcement members of uniform thickness with both horizontal and vertical components.\nHorizontal reinforcement and spacing is further taught in J. Pignataro in U.S. Pat. No. 6,629,393 issued Oct. 7, 2003, entitled “Masonry reinforcing tie.” Pignataro describes a ladder-type device formed of wire or rod with the parallel elongate members containing integrally formed, spacing elements. The patent describes providing lateral reinforcement and teaches only consistent horizontal mortar joint thickness in a masonry wall. Pignataro further teaches the use of roller dies to form the spacing elements.\nAccordingly, while several distinct devices were developed to assist in properly constructing a masonry wall unit, the current state of the art does not fulfill the need for a single efficient and economical combination device that simultaneously reinforces and spaces brick, block, stone or similar masonry building materials. As described hereinbelow, the present invention utilizes a combination reinforcing spacer device to horizontally reinforce and dually space both horizontally and vertically, thereby providing a useful and novel solution to the aforementioned difficulties."}
{"text": "The present invention relates generally to position determination, and particularly to position determination using DTV signals.\nThere have long been methods of two-dimensional latitude/longitude position location systems using radio signals. In wide usage have been terrestrial systems such as Loran C and Omega, and a satellite-based system known as Transit. Another satellite-based system enjoying increased popularity is the Global Positioning System (GPS).\nInitially devised in 1974, GPS is widely used for position location, navigation, survey, and time transfer. The GPS system is based on a constellation of 24 on-orbit satellites in sub-synchronous 12 hour orbits. Each satellite carries a precision clock and transmits a pseudo-noise signal, which can be precisely tracked to determine pseudo-range. By tracking 4 or more satellites, one can determine precise position in three dimensions in real time, world-wide. More details are provided in B. W. Parkinson and J. J. Spilker, Jr., Global Positioning System-Theory and Applications, Volumes I and II, AIAA, Washington, D.C. 1996.\nGPS has revolutionized the technology of navigation and position location. However in some situations, GPS is less effective. Because the GPS signals are transmitted at relatively low power levels (less than 100 watts) and over great distances, the received signal strength is relatively weak (on the order of −160 dBw as received by an omni-directional antenna). Thus the signal is marginally useful or not useful at all in the presence of blockage or inside a building.\nThere has even been a proposed system using conventional analog National Television System Committee (NTSC) television signals to determine position. This proposal is found in a U.S. patent entitled “Location Determination System And Method Using Television Broadcast Signals,” U.S. Pat. No. 5,510,801, issued Apr. 23, 1996. However, the present analog TV signal contains horizontal and vertical synchronization pulses intended for relatively crude synchronization of the TV set sweep circuitry. Further, in 2006 the Federal Communication Commission (FCC) will consider turning off NTSC transmitters and reassigning that valuable spectrum so that it can be auctioned for other purposes deemed more valuable."}
{"text": "This disclosure relates generally to the field of optics and, more specifically, to a method and apparatus for synthesizing ultra-wide bandwidth optical waveforms.\nAs waveforms are synthesized with increasing larger bandwidths, radio frequency (RF) design tends to grow in complexity. RF component design criteria such as power, amplitude ripple, gain flatness, phase distortion over bandwidth become increasingly more difficult to maintain with increasing bandwidths. What is need is an improved mechanism to increase waveform bandwidth of a radio frequency waveform during double sideband suppressed carrier modulation."}
{"text": "The present disclosure relates to non-transitory computer readable recording mediums that record a job transmission program for transmitting a job to an electronic device, and electronic devices that execute a job.\nElectronic devices of some type are capable of state transition between a job executable state and an energy saving state. The job executable state is a state in which a job is executable. While, the energy saving state is a state in which power consumption is less than that in the job executable state. In another electronic device, when no job is executed for a specific time period in the job executable state, the state is changed to the energy saving state."}
{"text": "Complementary Metal Oxide Semiconductor (CMOS) devices are dominated by n-channel (NMOS) and p-channel (PMOS) transistor structures. Various physical characteristics of each type of transistor determine the threshold voltage (Vt) that must be overcome to invert the channel region and cause a given transistor to conduct majority carriers (either by electrons movement in an NMOS device or by hole movement in a PMOS device).\nOne of the controlling physical characteristics is the work function of the material used to form the gate electrode of the transistor device. In semiconductor devices, such as a Dynamic Random Access Memory (DRAM) device, the transistor gates are predominantly made of polysilicon and an overlying layer of metal silicide, such as tungsten silicide and the gate dielectric is typically a high quality silicon oxide. The industry has moved to a transistor gate dielectric possessing a high dielectric constant of seven or greater (high-k dielectrics) for better leakage at given Effective Oxide Thickness (EOT). However, choosing a material with the appropriate work function as a gate electrode is still a challenge.\nStudies have been conducted in one area of using high-k dielectric concerning Fermi-level pinning at the polysilicon/metal oxide interface of the transistor gate structure. Taking HfO2, for example, the hafnium and the polysilicon form Hafnium-Silicon bonds whose energy level in the band gap causes the Fermi-Level of the polysilicon to be pinned near the conduction band. With this scenario, using the HfO2 as the transistor gate dielectric in an NMOS device, a small shift in the transistor Vt, relative to N+ polysilicon on SiO2 will occur due to Hafnium-Silicon interface. However, applying this case in a PMOS device, a large shift in the transistor Vt will occur due to the Hafnium-Silicon bonds still pinning the Fermi-Level of the polysilicon near the conduction band.\nTaking Al2O3, for example, the aluminum and the polysilicon form Aluminum-Silicon bonds that cause the Fermi-Level of the polysilicon to be pinned near the valence band due to the creation of the interface states that reside close to the valence band. With this scenario, using the Al2O3 as the transistor gate dielectric in a PMOS device, a small shift in the transistor Vt will occur due to P+ Aluminum-Silicon interface. However, applying this case in an NMOS device a large shift in the transistor Vt will occur due to the Aluminum-Silicon interface still having the Fermi-Level of the polysilicon being pinned near the conduction band and the transistor will not function in the desired range.\nCMOS transistor devices that use the traditional polysilicon gate electrodes in combination with a metal oxide dielectric (high-k dielectric) must be fabricated such that the NMOS and PMOS devices will each possess a suitable transistor threshold voltage (Vt).\nThere is a need for the construction of CMOS devices using high-k dielectric materials for the transistor gate dielectric which will successfully be used to form both n-channel (NMOS) and p-channel (PMOS) transistors in semiconductor devices."}
{"text": "A. Field of the Invention\nThe present invention relates to the field of insect screen assemblies, more specifically, insect screen assemblies that attach via a zipper, a hook and loop fastener, a plurality of snap buttons, or a hook and grommet configuration.\nB. Discussion of the Prior Art\nThe Fidler, Jr. Patent (U.S. Pat. No. 4,569,362) discloses a screened tent having panels that are detachable from the frame. However, the device disclosed under the Fidler Patent does not utilize a zipper as the means to connect the detachable panels.\nThe Valles Patent Application Publication (U.S. Pub. No. 2005/0055964) discloses a screened room enclosure that has detachable panels. However, the panels disclosed under the Valles Publication do not use a zipper as an attaching means to the rest of the enclosure.\nThe Benedyk Patent (U.S. Pat. No. 5,046,546) discloses a screened enclosure structure in which the screen panels are detachably removable from the frame member. However, the screen panels disclosed under the Benedyk Patent are not attached by a zipper.\nThe Zadok Patent (U.S. Pat. No. 6,192,643) discloses a modular screened enclosure system which is used for a pool deck or the like. However, the enclosure system disclosed under the Zadok Patent does not include detachable screen panels that are attached via a zipper attaching means.\nThe Narron Patent (U.S. Pat. No. 6,032,684) discloses a privacy screen for use on a deck or porch which has detachable panels. However, the detachable panels disclosed under the Narron Patent do not use a zipper as n attaching means.\nThe Antinone Patent (U.S. Pat. No. 3,763,917) discloses a detachable screen which is used in protecting porches, terraces, and the like. However, the detachable screen does not utilize a zipper as the attaching means.\nThe Hansen Patent (U.S. Pat. No. Des. 461,562) illustrates a design for a screened enclosure.\nIn light of the above discussed prior art there is a need for a detachable screen assembly involving a zipper, hook and grommet, nylon hook and loop, or snap buttons as the attaching means."}
{"text": "1. Field of the Invention\nThe present invention relates to a golf ball for a putting practice, and in particular to an improved golf ball for a putting practice in which a golf ball rolls and is moved in a desired direction by forming a concave groove or convex belt on a spherical surface of a golf ball.\n2. Description of the Background Art\nAs shown in FIG. 2, a groove represents a concave portion formed in a spherical surface of a golf ball and having a certain width and depth, and as shown in FIG. 12 a belt represents a convex portion formed a spherical surface of the same and having a certain width and height.\nResearches have been performed for many years for changing a size, depth and arrangement of the dimples formed on a surface of a golf ball for implementing a desired hydrodynamic property.\nFIG. 1 is a perspective view illustrating a conventional golf ball. As shown in FIG. 1, the golf ball 1 has a plurality of dimples formed on a surface 2 of the same at a certain interval therebetween for thereby implementing a desired hydrodynamic property. When a golf practice player puts a golf ball, the golf ball rolls based on a putting force and putting angle and a state of a practice field and then is moved into a hole. The putting practice is directed to straightly rolling a golf ball in a direction of a hole. The golf ball may be rolled and moved in various directions based on the condition of the practice field. Therefore, in order to test and enhance a player's ability, the golf ball must be rolled and moved in a desired direction when practicing the golf ball. At this time, the players practice the putting using a golf ball which is generally used in an actual golf field.\nHowever, in the conventional art, it is impossible for a practice player to select a certain golf ball based on his ability and to select the same based on his advanced ability."}
{"text": "Power hand tools, such as motorized ratchet wrenches and drivers, are commonly used in automotive, industrial and household applications to install and remove threaded fasteners and to apply a torque and/or angular displacement to a work piece such as a threaded fastener, for example. Motorized hand tools, such as cordless power ratchets and drivers, generally include an electric motor contained in a clamshell housing along with other components such as switches, light emitting diodes (LEDs), and batteries, for example. The clamshell housing generally includes two or more housing portions fastened together by fasteners such as screws or rivets. Windows in the clamshell housings may be provided for mounting the LEDs. The windows may be overmolded with a transparent or semitransparent material such as rubber or plastic to allow passage of light from the LEDs while protecting the internal components from water, dust and other foreign material.\nSecurely locating an LED in the clamshell housing has previously been accomplished by mounting the LED to a separate component such as a snap fit plastic insert or printed circuit board (PCB) before installing the LED and the separate component in the housing. These techniques have increased manufacturing costs by increasing the assembly steps and number of parts used to manufacture the motorized hand tool."}
{"text": "The present invention relates to semiconductor storage devices, and in particular, to a semiconductor storage device having a nonvolatile memory for executing a collective erase verify operation.\nNonvolatile semiconductor storage devices have been remarkably developed in recent years, and in particular, the rewriting times of flash memories have been substantially shortened. However, about a half of the rewriting time is consumed by the verify operation of checking whether or not the data in a memory cell has been rewritten in conformity to an expected value, and therefore, it is indispensable to provide a device for reducing this verify operation time.\nIn view of the above, there has been an attempt at reducing the verify operation time by providing sense amplifiers corresponding in number to memory cells to be subjected to a simultaneous write operation and simultaneously verifying the memory cells subjected to the simultaneous write operation in the verify operation.\nAccording to the above verify method, in regard to the write operation of the above rewrite operation, the number of memory cells to be subjected to the simultaneous write operation correspond to about 1 bit to 4 kbits. Since the data to be written into individual memory cells differ from one another, it can be considered efficient to simultaneously verify the memory cells of about 1 bit to 4 kbits subjected to the simultaneous write operation.\nHowever, in regard to the erase operation of the rewrite operation, memory cells of not less than 2 kbits to 512 kbits are collectively erased and all the memory cells are rewritten into identical data (\"0\" or \"1\"). Therefore, the conventional verify operation, in which memory cells of about 1 bit to 4 kbits are simultaneously verified by sense amplifiers similar to the write operation, is hard to be considered efficient.\nIn view of the above, a variety of erase verify methods as follows have been proposed. It is to be noted that the term \"verify\" mentioned hereinafter means the erase verify.\nThe discrimination between 0 and 1 of information held in the memory cells of the nonvolatile semiconductor storage device is normally executed by detecting whether a threshold voltage of the memory cell transistor is high or low. However, the verify methods described below are the verify methods of a memory cell array in which the low threshold voltage state is defined as an erased state. The verify operation in this case is required to collectively check the event that the threshold voltages of all the memory cell transistors have been reduced (i.e., the event that a current flows through all the memory cell transistors when a voltage intermediate between the high and low threshold voltages is applied to the gates of all the memory cell transistors).\n(1) The erase operation is completed at the point of time when the event that a current has flowed through the memory cell transistor exhibiting the minimum threshold voltage among all the memory cell transistors is detected (Japanese Patent Laid-Open Publication No. HEI 4-3395.)\n(2) The n(n: integer) memory cells connected to one word line are simultaneously verified by the same number of decision circuits (Japanese Patent Laid-Open Publication No. HEI 8-227590). This verify method has the same concept as that of the line test proposed in connection with the prior art DRAM (Dynamic Random Access Memory) and the like.\n(3) In a virtual ground type memory cell array, when a current flows by applying a voltage in series across the sources and the drains of a number of memory cell transistors connected to one word line, it is decided that all the memory cells connected to the word line have been erased (Japanese Patent Laid-Open Publication No. HEI 7-111901).\nHowever, the aforementioned prior art nonvolatile semiconductor storage device verify methods have the problems as follows.\nThat is, as stated before, the above verify methods are executed by collectively checking whether or not a current flows through all the memory cell transistors when the voltage intermediate between the high and low threshold voltages is applied to the gates of all the memory cell transistors.\nHowever, the operation of collectively verifying the event that a current flows through all the memory cells is required to detect one memory cell through which no current flows among a number of memory cells through which currents are flowing, and this is very hard to be achieved physically.\nFor example, according to the verify method listed in the item (1), the erase operation is executed by reducing the threshold voltages of the memory cell transistors from a programmed state D1 to an erased state D2 as indicated by the memory cell transistor threshold voltage distribution shown in FIG. 13. Then, the completion of the verify operation is determined by detecting the event that the threshold voltage A of the memory cell which is most likely to be erased has lowered from a word line selection voltage C in the verify operation. Accordingly, there still remains the possibility that the threshold voltage B of the memory cell which is least likely to be erased is greater than the selection voltage C and is kept intact in the programmed state, and this poses the problem that the erasing of all the memory cells has been verified.\nAccording to the verify method listed in the item (2), the n memory cells connected to one word line are simultaneously verified by the same number of decision circuits, and therefore, the problem of the verify method (1) can be solved. However, the method has the problem of an increase in area for the requirement of n decision circuits identical in number to the n memory cells to be verified collectively and the problem of an increase in time of the verify operation for the requirement of the verify operation times in a cycle identical to the number m(m: integer) of the word lines.\nAccording to the verify method listed in the item (3), the problem of the verify method (1) can be solved. However, if the ON-state resistance of the memory cell transistor and the substrate effect are taken into consideration, the method has the problem that the current is too infinitesimal and the problem that not so much memory cell transistors can be connected and simultaneously verified due to an increase in threshold value."}
{"text": "One example of a so-called optical interconnection structure for connecting semiconductor chips, primarily high-speed LSI circuits, with an optical signal is shown in FIGS. 1A and 1B. In this structure, LSI circuits 2 on substrate 1 are connected to each other by thin optical waveguide 5 parallel to substrate 1 through PD (PhotoDiode: light-detecting device) 3 and VCSEL (Vertical Cavity Surface-Emitting Laser) 4. In the example shown in FIGS. 1A and 1B, interposers 1a are provided. The number of general connections (parallelism) is in the range from about 100 to 1000. An optical circuit generally refers to an optical fiber, an optical waveguide, and a combination of an optical fiber and an optical waveguide."}
{"text": "The process of the present invention relates to the manufacture of personalized items such as jewelry. More particularly, the process of the present invention relates to an automated system that receives custom orders for personalized rings (i.e., class, championship, and affiliation) and generates the machining instructions that enable a milling machine to create the personalized ring from a wax blank.\nClass rings have been a popular keepsake among students for generations. Originally, they were relatively uniform and provided students little opportunity to express themselves. Over time, automated manufacturing processes made it possible to provide students customizing choices. Modern students are driving the class ring market toward a level of customization that has been previously economically impractical using present manufacturing methods.\nPresent manufacturing methods include the use of computer aided design/computer-aided manufacturing (CAD/CAM). CAD/CAM has facilitated producing customized rings in large quantities. The present level of customization provides personalized features such as: student's name, school name, graduation year, icons, academic degrees, and the like.\nTraditionally, the use of CAD/CAM in the jewelry industry has been primarily focused on the manufacture of custom molds and engraving or otherwise machining the jewelry directly. These two approaches have limitations. Machining molds using CAD/CAM is too expensive for single-use custom applications. Engraving jewelry is also expensive due to the precious metal lost to scrap, manufacturing errors and ordering errors.\nCAD/CAM technology is also difficult to automate for the purpose of making personalized products. In one legacy system, a CAD/CAM operator manually manipulates a geometric model of a ring by grabbing a surface on the blank geometric model, defining the boundary splines, projecting the text or graphic onto the surface and then instructing the CAD/CAM software to generate machining instructions for the geometric model that has been created. The machining instructions result in a desired toolpath for a computer numerically controlled (“CNC”) milling machine. Using human operators to repeat these steps manually in order to generate the machining instructions for thousands of individual, personalized rings is cost prohibitive."}
{"text": "An electronic control unit in which high-heat-generating devices, such as switching devices, are mounted on a circuit board is known. The high-heat-generating device generates a lot of heat when the electronic control unit operates. For example, in an electronic control unit disclosed in JP-A-2011-23459 corresponding to US 2011/0013370A1, switching devices as high-heat-generating devices are mounted on a surface of a circuit board, and a surface of a heat-dissipating member is located facing the surface of the board.\nIn the electronic control unit disclosed in JP-A-2011-23459, the switching device is held in a room formed in the heat-dissipating member. However, the room is not shaped to match the shape of the switching device, and the capacity of the room is a few times larger than the volume of the switching device. For this reason, in the assembled condition of the electronic control unit, there is a large distance between the inner wall of the room and the outer wall of the switching device. Accordingly, it may be difficult to efficiently transfer heat of the switching device to the heat-dissipating member. As a result, the heat of the switching device is not efficiently dissipated through the heat-dissipating member.\nFurther, in the assembled condition of the electronic control unit disclosed in JP-A-2011-23459, there is a large distance between the facing surfaces of the heat-dissipating member and the board. Accordingly, it may be difficult to transfer the heat of the switching device to the heat-dissipating member through the board.\nFurthermore, in the assembled condition of the electronic control unit disclosed in JP-A-2011-23459, a heat-conducting member is provided in each room. Accordingly, the number of man-hours needed to place the heat-conducting member may be increased.\nIn particular, the electronic control unit disclosed in JP-A-2011-23459 is used to control a motor of an electric power steering apparatus for a vehicle. When the motor of the electric power steering apparatus is driven and produces assisting-torque to help a driver to steer the vehicle, a large current flows through the motor and the high-heat-generating devices of the electronic control unit. Therefore, when the motor is driven, the high-heat-generating devices generate a lot of heat. For this reason, when an electronic control unit is used to control a motor of an electric power steering apparatus, it is preferable that heat of high-heat-generating devices of the electronic control unit should be dissipated efficiently."}
{"text": "In recent years the use of separate wheeled luggage carriers which are collapsible and having a telescoping handle assembly has been paralleled by a growth in the use of wheeled luggage in which, generally, four wheels were provided on the bottom of a rectangular parallelepipedal luggage article which could have a tether or handle at the top adapted to be pulled by the user to tow the article along the ground or floor.\nMore recently it has been found to be advantageous to provide the handle so that it can pull out form the luggage article directly.\nSuch systems have not been fully successful heretofore. For example, the pull-out handle assembly may not be fully reliable or may be subject to wear because of continuous friction between the handle and its guides.\nNot all luggage articles can adequately be provided with prior art handle structures and, for example, catalog cases cannot adequately be formed with conventional pull-out handle structures."}
{"text": "Three-dimensional imaging on a computer display has become more prevalent in recent years with the advancement of computing power and computer modeling. Active research and development in biological science and engineering have been ongoing for processing three-dimensional images for instructional purposes and medical applications. One area of particular interest stems from the medical imaging field for the treatments of patients to compare imaging scans of a patient taken before, during and after the treatments to assess the progress and effectiveness of the therapy or the degree of disease progression. Suitable medical image scanners to capture the image scans include x-ray, computerized tomography (CT), cone beam computerized tomography (CBCT), magnetic resonance imaging (MRI) and positron emission tomography (PET).\nMedical imaging technology measured by both qualitative and quantitative analyses has continued to be developed with the availability of software applications and computing capabilities. Technical challenges in combining hundreds of sectional images of CT scan to produce a three-dimensional (“3D”) image representation of a body organ have been hindered by some principal issues, such as conversion methodologies, surface color manipulation and insufficient image resolution.\nIt is therefore desirable to have methods and systems for qualitative and quantitative transformations on multiple datasets of sectional images into overlapping and delineated three-dimensional image representations."}
{"text": "The present invention relates to medical instruments for ablating tissue, and in particular to a microwave probe for ablation of tumors and the like.\nMicrowave ablation (MWA), like radio frequency ablation (RFA), uses localized heating to cause tissue necrosis. However, MWA can produce greater and more rapid heating and can easily support the use of multiple probes because current flow between the probes can be limited. The mode of heating in MWA also eliminates ground pads and charring concerns.\nUnfortunately, current MFA equipment produces relatively small lesions because of practical limits in power and treatment time. Power is limited by the current carrying capacity of the small gauge feeder line as it passes through the patient to the site of the necrosis. Larger feeder lines are undesirable because they are not easily inserted percutaneously. Heating of the feeder line at high powers can also lead to burns around the insertion point of the MWA probe."}
{"text": "1. Technical Field\nThis document relates to materials and methods for circulatory bypass of a ventricle in the heart of a mammal. For example, materials and methods for bypassing a permanently or temporarily impaired left ventricle in the heart of a mammal (e.g., a human) are provided.\n2. Background Information\nCongestive heart failure (CHF) affects the ability of the heart to provide sufficient blood flow to the organs of the body. Conditions such as coronary artery disease, scar tissue from a myocardial infarction, high blood pressure, and heart valve disease can contribute to congestive heart failure. Pharmaceutical intervention can expand blood vessels and/or or lower blood pressure to allow blood to flow more easily and the heart to pump more efficiently. If the damage to the heart is extensive, surgical intervention can be required."}
{"text": "(1) Field of the Invention\nThe present invention relates to a structure of a water-cooled engine, in particular relates to a snowmobile engine cooling arrangement having cooling water passages for cooling a snowmobile engine.\n(2) Description of the Prior Art\nConventionally, most small snow vehicle such as snowmobiles and the like use two-cycle engines because they are relatively simple in structure, light and compact and yet powerful. Recently, however, because of regulation of exhaust gas due to environmental issues and aiming at improvement in reduction of fuel consumption rate, there is a trend toward employing four-cycle engines.\nAs the engine structure for snowmobiles, water-cooled engines have been generally employed because they can stably cool the engine by cooling water, preventing the occurrence of overheating or overcooling and improving the output power and yet are effective in reducing noise.\nAs an example of conventional engine cooling arrangements, Japanese Patent Application Laid-Open No. 2001-12243 discloses a cooling system for a snowmobile engine. This cooling system is comprised of a cooling water jacket formed inside the engine, a heat exchanger for cooling the cooling water and a cooling water pump for transferring the cooling water so as to circulate the cooling water inside the engine, and further includes: a cooling water bypass through which the cooling water can pass, avoiding its flow into the heat exchanger; and a switching valve for permitting the cooling water to flow into the cooling water bypass, wherein the switching valve is adapted to switch the cooling water circuit in such a manner that the cooling water is allowed to pass through the heat exchanger when the cooling water is equal to or higher in temperature than a predetermined level, and the cooling water is restricted from passing through the heat exchanger and is conducted to the cooling water bypass when the temperature of the cooling water is lower than the predetermined temperature level.\nMeanwhile, in contrast to two-cycle engines, which are high in power, light and compact, four-cycle engines need a camshaft and oil lubrication, inevitably tending towards large size.\nIn particular, when a four-cycle engine is included in a snowmobile, it is necessary to provide a contrived layout of the oil pan configuration, intake and exhaust systems and associated auxiliaries, in order to make the body and engine hood have similar size to those of a two-cycle engine.\nFurther, in general, when the mounted engine is of a water cooled type, a reservoir tank should be interposed within the piping of the cooling water. That is, it is necessary to take into account the installation space for this reservoir tank. Further, the reservoir tank needs to be arranged at a site where it is easy for cooling water to be re-supplied.\nSince the snowmobile is used in cold areas, in order to stabilize the idling engine speed at the start of operation or in order to regulate the surging of the thermostat, it is necessary to consider the layout of the cooling water piping and the position of the thermostat.\nFurther, the cooling water pump for transferring cooling water is driven by a belt which is driven by the rotation of the crankshaft. This means the enlargement of the longitudinal size of the engine. Therefore it is necessary to arrange the cooling water pump at an appropriate position. Moreover, it is necessary to choose the position of the cooling water pump optimally, considering the positional relationship with an alternator, which is belt driven, so that the front-to-rear size of the engine will be minimized.\nThe present invention has been devised in view of the above conventional problems, it is therefore an object of the present invention to provide a cooling arrangement for a snowmobile engine, which is able to stabilize the engine idling speed at the start of operation and is improved in the work performance in supplying cooling water and realizes a space-saving engine layout.\nIn order to achieve the above object, the present invention is configured as follows:\nIn accordance with the first aspect of the present invention, a cooling arrangement for a snowmobile engine, includes: a cooling water jacket formed inside an engine; a heat exchanger for cooling water; a water pump for ejecting cooling water; and a thermostat for controlling supply of the cooling water ejected from the water pump into the engine, so as to cool the engine by circulating the cooling water inside the engine, and is characterized in that:\nthe engine is mounted with its cylinder head at top, and an exhaust manifold is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder head while an intake manifold is disposed on the rear side, with respect to the vehicle\"\"s direction of travel, of the cylinder head;\nthe water pump is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of cylinder block of the engine, under the exhaust manifold and is coupled to a crankshaft projected from one side wall, with respect to the widthwise direction of the body, of the cylinder block, so that it is driven by the rotational force which is transmitted through a drive belt from a rotational member fitted on one end of the crankshaft;\na cooling water inlet port for leading cooling water into the cooling water jacket inside the engine is provided on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and under the exhaust manifold at a position close to the drive belt;\na cooling water outlet port through which cooling water is taken out from the interior of the engine is arranged on the engine\"\"s side face at a position over and opposite to the drive belt; and\nthe thermostat is arranged at the cooling water outlet port, at a point downstream with respect to the flow of cooling water.\nIn accordance with the second aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the cooling water outlet port is comprised of a thermo housing for incorporating the thermostat and a thermo cap for covering the thermostat, the thermo cap being connected to the heat exchanger by way of a cooling water passage.\nIn accordance with the third aspect of the present invention, the cooling arrangement for a snowmobile engine having the above second feature is characterized in that the thermo housing is formed with a first cooling water bypass passage which branches off at a point upstream, with respect to the flow of cooling water, of the thermostat and is connected by way of a throttle body to the cooling water passage at a point upstream of the water pump.\nIn accordance with the fourth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above second feature is characterized in that the thermo housing is formed with a second cooling water bypass passage which branches off at a point upstream of the thermostat and is connected directly to the cooling water passage at a point upstream of the water pump without passing through the throttle body.\nIn accordance with the fifth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above third feature is characterized in that the thermostat is used to control the outlet of cooling water after passage of the cooling water jacket inside the engine while the thermo housing side and cooling water passage side are kept so as to be in constant communication to each other through the first cooling water bypass passage and the second cooling water bypass passage.\nIn accordance with the sixth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above fourth feature is characterized in that the thermostat is used to control the outlet of cooling water after passage of the cooling water jacket inside the engine while the thermo housing side and cooling water passage side are kept so as to be in constant communication to each other through the first cooling water bypass passage and the second cooling water bypass passage.\nIn accordance with the seventh aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that, in the snowmobile, an oil filter is arranged on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and an oil cooling means is interposed between the oil filter and the cylinder block.\nIn accordance with the eighth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that, in the snowmobile, an oil cooling means is provided on the front side, with respect to the vehicle\"\"s direction of travel, of a crawler for causing the snowmobile to move.\nIn accordance with the ninth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that, in the snowmobile, a muffler is disposed in front of the engine body in the engine room and an oil cooling means is arranged in front of the muffler.\nIn accordance with the tenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the water pump is arranged between the engine and the exhaust passage provided in front of the engine.\nIn accordance with the eleventh aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the water pump is arranged over an oil pan.\nIn accordance with the twelfth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the water pump is arranged under the exhaust manifold.\nIn accordance with the thirteenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that the engine has an alternator provided at a position opposite the water pump with the cylinder block in between, the three components being laid from the front to rear with respect to the vehicle\"\"s direction of travel.\nIn accordance with the fourteenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that a cooling water reservoir tank is arranged in the rear of the engine, at the side of the intake manifold, over an oil pan, heat exchanger and alternator.\nIn accordance with the fifteenth aspect of the present invention, the cooling arrangement for a snowmobile engine having the above first feature is characterized in that a cooling water reservoir tank is arranged without being projected outwards beyond a drive belt for driving an alternator and water pump, when viewed from top.\nAccording to the present invention, the following effects can be obtained.\nFirst, in a cooling arrangement for a snowmobile engine, the engine is mounted with its cylinder head at top, and an exhaust manifold is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder head while an intake manifold is disposed on the rear side, with respect to the vehicle\"\"s direction of travel, of the cylinder head. The water pump is disposed on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder head of the engine, under the exhaust manifold and is coupled to a crankshaft projected from one side wall, with respect to the widthwise direction of the body, of the cylinder block, so that it is driven by the rotational force which is transmitted through a drive belt from a rotational member fitted on one end of the crankshaft. Further, a cooling water intake port for leading cooling water into the cooling water jacket inside the engine is provided on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and under the exhaust manifold, at a position close to the drive belt while a cooling water outlet port through which cooling water is taken out from the interior of the engine is arranged on the engine\"\"s side face opposite to the drive belt side or on the opposite side at a position above the crankshaft from which engine power is taken. This arrangement makes it possible for the cooling water to flow through the engine interior, approximately diagonally from the bottom to the top, hence the water is able to efficiently flow and be discharged without stagnation inside the cylinder head and inside the cylinder block. Accordingly, this configuration provides an engine of an improved cooling efficiency compared to the conventional configuration. Further, since supply and discharge of cooling water is achieved utilizing the dead space near the water pump and under the exhaust manifold, this arrangement realizes an engine of a highly efficient space usage.\nSince the cooling water outlet port is formed with a thermo housing incorporating a thermostat and a thermo cap covering the thermostat while the thermo cap is connected to the heat exchanger by way of a cooling water passage, e.g., a cooling water hose so as to control the cooling water at the cooling water outlet, by making a so-called outlet control, it is possible to perform correct cooling water control at a stable cooling water temperature.\nSince the thermo housing is formed with a first cooling water bypass passage which branches off at a point upstream, with respect to the flow of cooling water, of the thermostat and is connected by way of a throttle body to the cooling water passage at a point upstream of the water pump, it is possible to prevent icing and make stable intake control by maintaining the cooling water supplied to the throttle valve at an approximately fixed temperature.\nSince the thermo housing is formed with a second cooling water bypass passage which branches off at a point upstream, with respect to the flow of cooling water, of the thermostat and is connected directly to the cooling water passage at a point upstream, with respect to the flow of cooling water, of the water pump without passing through the throttle body, the cooling water warmed by the engine is supplied to the engine. This makes it possible to regulate the temperature of the cooling water at an approximately constant level under the thermostat control of supply of the cooling water cooled through the heat exchanger, hence realizing a good operational condition of the engine.\nSince control of cooling water by the thermostat is performed by regulating the outlet of cooling water after passage of the cooling water jacket inside the engine, it is possible to make reliable cooling water control based on the cooling water of a stabilized water temperature. Further, since the thermo housing side and cooling water passage side are kept so as to be in constant communication to each other through the first cooling water bypass passage and the second cooling water bypass passage, this makes it possible to regulate the temperature of the cooling water in the engine at an approximately constant level. As a result, the thermostat for opening and closing the main cooling water passage in response to the cooling water temperature quickly reacts to the cooling water temperature that varies dependent on the operational state of the engine, to thereby prevent the occurrence of engine seizure and other troubles.\nSince, in the snowmobile, an oil filter is arranged on the front side, with respect to the vehicle\"\"s direction of travel, of the cylinder block and an oil cooling means, e.g., an oil cooler, is interposed between the oil filter and the cylinder block, it is possible to efficiently cool not only the engine but also the engine oil, thus achieving a beneficial operational condition of the engine.\nSince, in the snowmobile, an oil cooling means is provided on the front side, with respect to the vehicle\"\"s direction of travel, of a crawler for causing the snowmobile to move, it is possible to perform improved cooling by cooling oil with the scattered snow. Further, since two divided heat exchanging means or heat exchangers are provided, where one is used as usual for cooling water and the other for cooling oil, it is possible to realize different cooling functions with a simple structure.\nSince a muffler is disposed in front of the engine body in the engine room in the snowmobile and an oil cooling means is arranged in front of the muffler, this arrangement allows the oil cooling means to receive the flow of air, during travel, ahead of the muffler and exhaust pipe which are higher in temperature than the oil cooler. Thus, it is possible to cool the oil at an improved efficiency.\nSince the water pump is arranged between the engine and the exhaust passage provided in front of the engine, it is possible to realize a space-saving engine layout utilizing the dead space under the exhaust passage.\nSince the water pump is arranged over an oil pan, it is possible to prevent the engine height from increasing.\nSince the water pump is arranged under the exhaust manifold, it is possible to realize a space-saving engine layout utilizing the dead space under the exhaust manifold.\nSince the engine has an alternator provided at a position opposite the water pump with the cylinder block in between, the three components being laid from the front to rear with respect to the vehicle\"\"s direction of travel, it is possible to realize a space-saving and well-balanced drive belt layout.\nSince a cooling water reservoir tank is arranged in the rear of the engine, at the side of the intake manifold, over an oil pan, heat exchanger and alternator, this arrangement of the reservoir tank at the top of the engine room makes it easy for cooling water to be re-supplied and other maintenance.\nSince a cooling water reservoir tank is arranged without being projected outwards beyond a drive belt for driving an alternator and water pump, when viewed from top, this arrangement realizes a compact engine configuration."}
{"text": "1. Field of the Invention\nThis invention relates to tyres for vehicle wheels, in particular tyres intended for use on the steering axle of trucks or articulated lorries for medium/heavy transport on long journeys, generally on motorways, at relatively high speed.\n2. Description of the Related Art\nMedium/heavy transport vehicles generally have different types of tyre fitted on the driving axles and steering axles depending on the different behavioural and functional prerogatives required of the tyres under the different working conditions in which they are used.\nIn general, tyres intended for use on the steering axle and for long motorway journeys need: good evenness of wear, combined with a good mileage; good directional stability and steering precision on dry surfaces; good grip on wet surfaces; good resistance with respect to tearing of the tread and to lateral shocks; low rolling resistance; and the ability to expel small stones trapped in the grooves. Factors contributing to producing these features are the tread design, the profile of the inflated tyre and the structure of the tyre.\nIn connection with the tread design, it is preferably of the type comprising a plurality of continuous circumferential ribs, two shoulder ribs and one or more centrally positioned ribs, delimited by a corresponding plurality of circumferential grooves.\nThe problem of uneven wear in tyres of this type is well-known in the art. For most of its life, the tyre is used for long journeys on straight stretches, with minimal stresses due to maneuvering and/or bends. In long journeys, the failure to optimize the distribution of the tyre/ground contact pressures on the tyre's footprint may lead to continuous microsliding of parts of the tread on the asphalt, consequently inducing premature wear of these parts of the tread and, in any case, uneven wear.\nOne of the parts of the tread particularly sensitive to this unevenness of wear is the shoulder zone comprising both the axially external edge of the shoulder rib and also the zone straddling the shoulder groove, i.e. the groove delimiting the aforementioned shoulder rib in the position axially internal to the latter. Various solutions have been proposed to solve the problem of uneven wear in the shoulder zone.\nThis invention derives from the perception of the technical problem of uneven wear in the shoulder zone of the tyre, dependent on the radially external profile of the surface of the tread in the said zone.\nU.S. Pat. No. 4,480,671 describes a tyre for wheels of vehicles for heavy transport characterized in that the tread comprises in each of the two lateral zones a narrow circumferential groove which defines a lateral rib having a meridian profile substantially parallel and lowered with respect to the meridian profile of the geometrical envelope of the median zone of the tread. The width of the said lateral rib is at most equal to the maximal depth of the groove; the difference in level between the surface of the lateral ribs and that of the surface of the median zone of the tread is between 40% and 200% of the deflection of the tread under load and is such that, under normal working conditions, the said lateral ribs are in contact with the ground.\nThe patent application WO 92/2380 describes a tyre for wheels of vehicles for heavy transport comprising a radial carcass, a belt structure comprising at least two strips of different axial width and a tread band having in the lateral zones a circumferential groove of width preferably between 0.6 and 8 mm, delimiting a lateral rib the meridian profile of which is substantially parallel and lowered with respect to the meridian profile of the geometrical envelope of the median zone. The difference between the width of the tread band and the axial distance between the median axes of the lateral grooves is between 34 and 80 mm. This solution improves the separation resistance between the extremities of the belt strips and between the extremities of the entire belt structure and the carcass structure.\nThe patent application WO 98/26945 describes a tyre for wheels of vehicles for heavy transport comprising a tread band having a plurality of main ribs defined by circumferential grooves and a sacrifical rib in both the lateral zones of the tyre, separated from the main ribs by a narrow shoulder groove having a maximum width of 1.5 mm. The sacrifical rib has a width of between 2.5% and 12% of the width of the tread band and a constant radial lowering with respect to the transversal profile of the main ribs of between 0.5 and 2 mm, when the tyre is new.\nThe abstract of patent JP 05-77608 describes a radial tyre for wheels of heavy transport vehicles in which a pair of thin circumferential grooves is defined on the tread band. A first arc of a circle defines the tread profile between the equatorial plane and the thin groove and a second arc of a circle defines the profile of the tread between the thin groove and the shoulder edge, in order to radially lower the tread profile in the region around the thin circumferential groove, both towards the centre of the tread and towards the shoulder edge.\nU.S. Pat. No. 5,660,652 describes a tyre for use on the steering axle comprising a tread having a plurality of continuous circumferential grooves, a pair of thin grooves in the shoulder zone, a pair of radially lowered ribs local to the shoulder edges and a plurality of main ribs. The tread has a radially external surface having a curvature initially defined by a radius of curvature internal to the tyre, followed by a curvature defined by a radius of curvature external to the tyre. The main shoulder rib has a progressively increasing thickness starting from the axially internal, adjacent circumferential groove and reaching a maximum value close to the thin groove. Straddling the said circumferential groove adjacent to the main shoulder rib, the tyre/ground contact pressures equalize, whereas a greater contact pressure is obtained on the shoulder edge. This distribution of tyre/ground contact pressures is obtained thanks to the progressive increase in thickness of the tread from the point of contraflexure corresponding to the inversion of the curvature to the thin shoulder groove. The added rubber effectively provides material in the tread shoulder region, where the propensity to wear is high, and increases the contact pressure of the region adjacent to the thin groove."}
{"text": "Conventionally, in production of a semiconductor device, microprocessing by lithography using a resist composition has been carried out. The microprocessing is a processing method in which a thin film of a photoresist composition is formed on a semiconductor substrate such as a silicon wafer, irradiated with an active light such as an ultraviolet light through a mask pattern that has a pattern of a device, and developed, and the substrate is etched using the obtained photoresist pattern as a protective film to form fine irregularities corresponding to the pattern on a surface of the substrate. In recent years, an increase in degree of integration of the semiconductor device has advanced, and the wavelength of used active light shortens from an i line (wavelength: 365 nm) and a KrF excimer laser (wavelength: 248 nm) to an ArF excimer laser (wavelength: 193 nm). For this reason, diffuse reflection of active light on the semiconductor substrate and influences of standing wave are severe problems. In order to solve the problems, a method for providing an anti-reflective coating (Bottom Anti-Reflective Coating: BARC) between the resist and the semiconductor substrate has been widely investigated. The anti-reflective coating is also referred to as resist underlayer film. As the anti-reflective coating, an organic anti-reflective coating formed from a polymer having a light absorbing moiety and the like has been often investigated from the viewpoint of ease of use thereof.\nPatent Documents 1 to 3 disclose a resist underlayer film (anti-reflective coating) in which no intermixing with a photoresist film formed thereon is caused, a desired optical parameter (k value and n value) is obtained during exposure using an ArF excimer laser, and a desired dry etching rate is obtained.\nOn the other hand, in lithography using EUV (abbreviation of extreme ultraviolet light, wavelength: 13.5 nm) exposure that is a further microprocessing technology, reflection on a substrate does not occur, but roughness of side wall of a resist pattern due to miniaturization of a pattern becomes a problem. Therefore, a resist underlayer film for forming a resist pattern profile having high rectangularity has been often investigated. As a material for forming a resist underlayer film for exposure to a high-energy beam such as EUV, an X ray, and an electron beam, a resist underlayer film-forming composition in which generation of outgas is reduced has been disclosed. (see Patent Document 4)."}
{"text": "Conventionally, for example, Japanese Laid-Open Patent Publication No. 80791/1987 (Tokukaishou 62-80791) has disclosed a method for recognizing a character image that has been picked up, wherein: image information, such as a character image that has been inputted, is set as a reference pattern, the storage size of a dictionary pattern is set based on the size of the reference pattern, and the image information is recognized based on the dictionary pattern.\nMoreover, Japanese Laid-Open Patent Publication No. 74729/1994 (Tokukaihei 6-74729) has disclosed a method wherein: the distance between a character to be recognized and a camera is calculated, a zooming operation of the camera is carried out based on the calculated distance so as to allow the character to be recognized to have a predetermined size so that the character image to be recognized is inputted.\nHowever, the problem with the character recognition device using the conventional methods is that in the case when a character image, which has been temporarily stored, is recognized later, if the character image once stored can not be recognized, the image pick-up operation (reading operation) has to be repeatedly carried out until the character image has been recognized.\nFor this reason, for example, in the case when a character image, which has been picked up by using a portable-type electronic camera serving as a character recognition device and has been stored therein, is later subject to a character-image recognition at a site different from the site where it was picked up, if the stored character image can not be recognized, the user has to return to the site where it was picked up and re-pick up the same character image.\nFurthermore, at the time when a character image is picked up, neither of the character recognition devices using the above-mentioned conventional methods makes a judgment as to whether or not the character image can be recognized. For this reason, even in the case when the character image that has been picked up is recognized at that site, if the character image that has been picked up includes any character that can not be recognized, the picking-up operation (the reading operation) of the character image has to be again carried out upon completion of the recognizing operation.\nAnother problem with the above-mentioned conventional techniques is that the enlarging operation is applied to even recognizable character images unnecessarily."}
{"text": "1. Field of the Invention\nThe present invention relates to the transmission of data, and in particular to the transmission of digital data coded by means of a turbocoder.\n2. Discussion of the Related Art\nA turbocoder includes two (or more than two) simple convolutional type coding circuits and one or several inter-leavers operating on large data blocks (typically, from 100 to 16,000 bits).\nFIG. 1 illustrates a conventional four-state convolutional type coder. Coder 1 receives, at each pulse of a clock, not shown, a bit D. Bit D belongs to a bit sequence corresponding to data to be transmitted, digitized and arranged in series.\nData bit D is directly provided at the output of coder 1 by a first branch 11. Bit D is provided to a coding circuit C by a second branch 12. Circuit C outputs a so-called redundancy bit Y.\nCircuit C includes four D-type flip-flops Bi (B0, B1, B2, B3) connected in series and five adders 2i (20, 21, 22, 23, and 24). Adders 2i are adders modulo 2. At each clock pulse, the state (“0” or “1”) at the input of a flip-flop Bi is transferred to its output. Adders 2i are coupled to elements of coder 1 (input, flip-flops Bi, other adders, output) so that redundancy bit Y represents a linear combination of four successive data bits."}
{"text": "Estimating relatedness of various data sets is useful in many applications. In social networking, for example, estimating potential interest of a first user in another user may be useful when selecting users to suggest as potential “friends” to the first user. It is also useful in selecting advertising to present (render) to the particular user. Similarly, a virtual space can be defined so that it describes a given number of qualities. For example, a virtual reality space can be defined so that the distance between avatars represents similarity of their qualities rather than the distance in the “physical” space as it appears on a computer screen or other display. There are other uses for estimating relevance/relatedness as well.\nTherefore, it is desirable to facilitate the process of estimating relatedness of real or notional objects, and taking action based on the resulting estimates."}
{"text": "Aluminium profiles are used extensively in applications such as window and door frames, display counters and partitions. Such aluminium profiles come in several finishes, among which are metallic finishes such as the well-known silver-metallic and gold or bronze metallic finishes. Recently a brown-metallic finish has become popular as well. These finishes involve anodizing the aluminium. Alternatively, aluminum profiles have been coated with paints which are baked onto the aluminum. To anodize aluminum, it has to undergo an electro-chemical process which is rather complex, expensive and limits the colours obtainable. Aluminum profiles that are painted require first to undergo a preliminary process of undercoating in order to provide a suitable surface for the paint to adhere to, since paints do not adhere well to aluminum.\nAlthough many advances and developments have taken place in the coating industry, including the art of finishing aluminum profiles, there are a number of problems that have as yet not been overcome. First of all the finished aluminum profiles, which are used for so many diverse applications, require special resistance against degradation by light, chemicals, corrosive atmospheric gases and other pollutants, because they are used so much under outdoor conditions in displays and window frames where ultra-violet light, industrial gases and other weather conditions deteriorate the coating and the aluminum. This is particularly true for use of aluminum near coastal areas, where oxygen and salt water in the air subjects the profiles to severe corrosion.\nA further problem is caused by the fact that modern household and sanitary cleaning materials contain bleaches and alkali which react with aluminum and deteriorate its surface. These problems dictate that the aluminum profiles for use in such applications must have complete protection from contact between any corrosive atmosphere and the aluminum. This protection is very difficult to achieve in profiles with intricate channels and grooves, as is the case with aluminum profiles. Yet the demand for aluminum profiles for such applications is ever-increasing, particularly for profiles with metallic finishes at a reasonable price."}
{"text": "It is known that carbon-carbon double bonds in polymers may be successfully hydrogenated by treating the polymer in an organic solution with hydrogen in the presence of a catalyst. Such processes can be selective in the double bonds which are hydrogenated so that, for example, the double bonds in aromatic or naphthenic groups are not hydrogenated and double or triple bonds between carbon and other atoms such as nitrogen or oxygen are not affected. This field of art contains many examples of catalysts suitable for such hydrogenations, including catalysts based on cobalt, nickel, rhodium, ruthenium, osmium, and palladium. The suitability of the catalyst depends on the extent of hydrogenation required, the rate of the hydrogenation reaction and the presence or absence of other groups, such as carboxyl and nitrile groups, in the polymers. There have been already many patents and patent applications filed in this area, such as U.S. Pat. No. 6,410,657, 6,020,439, 5,705,571, 5,057,581, and 3,454,644.\nHowever, many diene-based polymers, copolymers or terpolymers are made by emulsion polymerization processes and they are in latex form when they are discharged from polymerization reactors. Therefore it is very desirable to directly hydrogenate a diene-based polymer in latex form and this approach is receiving increasing attention in the recent decade. Many efforts have been made to realize such a process.\nSo far significant attention has been paid to the hydrogenation of C═C bonds using hydrazine or a derivative of hydrazine as a reducing agent together with an oxidant like oxygen, air or hydrogenperoxide. The hydrogen source to saturate the C═C bonds is then generated in-situ as a result of the redox reactions in which diimide is also formed as intermediate.\nIn U.S. Pat. No. 4,452,950 the latex hydrogenation is performed using the hydrazine hydrate/hydrogen peroxide (or oxygen) redox system to produce diimide in situ. CuSO4 or FeSO4 is used as a catalyst.\nU.S. Pat. No. 5,039,737, and 5,442,009 provide a more refined latex hydrogenation process which treats the hydrogenated latex with ozone to break the cross-linked polymer chains which form during or after the latex hydrogenation using the diimide approach.\nU.S. Pat. No. 6,552,132 B2 discloses that a compound can be added before, during or after the latex hydrogenation to break crosslinks formed during the hydrogenation using the diimide hydrogenation route. The compound can be chosen from primary or secondary amines, hydroxylamine, imines, azines, hydrazones and oximes.\nU.S. Pat. No. 6,635,718 B2 describes the process for hydrogenating C═C bonds of an unsaturated polymer in the form of an aqueous dispersion by using hydrazine and an oxidizing compound in the presence of a metal compound containing a metal atom in an oxidation state of at least 4 (such as Ti(IV), V(V), Mo(VI) and W(VI)) as the catalyst.\nIn Applied Catalysis A-General, Vol. 276, no. 1-2, 2004, 123-128 and Journal of Applied Polymer Science, Vol. 96, no. 4, 2005, 1122-1125 detailed investigations relating to the hydrogenation of nitrile butadiene rubber latex via utilization of the diimide hydrogenation route are presented which cover examining hydrogenation efficiency and degree of hydrogenation. It has been found that there are side reactions at the interphase of the latex particles and within the polymer phase, which generate radicals to initiate the crosslinking of polymers in the latex form. Using radical scavengers did not show any evidence in helping to suppress the degree of gel formation.\nAlthough there are methods developed to reduce the crosslinking, the aforementioned diimide route still encounters gel formation problem, especially when high hydrogenation conversion is achieved. Therefore, the resulting hydrogenated rubber mass is difficult to process or is unsuitable for further use because of its macroscopic three dimensional cross-linked structure.\nU.S. Pat. No. 5,272,202 describes a process for the selective hydrogenation of the carbon-carbon double bonds of an unsaturated, nitrile-group-containing polymer with hydrogen in the presence of a hydrogenation catalyst being a palladium compound. In this process an aqueous emulsion of the unsaturated, nitrile-group-containing polymer is subjected to the hydrogenation and additionally an organic solvent capable of dissolving or swelling the polymer is used at a volume ratio of the aqueous emulsion to the organic solvent in a range of from 1:1 to 1:0.05. The aqueous emulsion is brought into contact with gaseous or dissolved hydrogen while maintaining an emulsified state.\nU.S. Pat. No. 6,403,727 discloses a process for selectively hydrogenating C═C double bonds in polymers. Said process involves reacting the polymers with hydrogen in the presence of at least one hydrogenation catalyst selected from the salts and complex compounds of rhodium and/or of ruthenium, in an aqueous dispersion of the polymers which comprises up to 20% by volume of an organic solvent. The suitable rhodium containing catalysts are rhodium phosphine complexes of the formula Rh XmL3L4(L5)n wherein X is a halide, the anion of a carboxylic acid, acetylacetonate, aryl- or alkylsulfonate, hydride or the diphenyltriazine anion and L3, L4 and L5 independently are CO, olefins, cycloolefins, dibenzophosphol, benzonitrile, PR3 or R2P-A-PR2, m is 1 or 2 and n is 0, 1 or 2, with the proviso that at least one of L3, L4 or L5 is one of the above mentioned phosphorus-containing ligands of the formula PR3 or R2-A-PR2, wherein R is alkyl, alkyloxy, cycloalkyl, cycloalkyloxy, aryl or aryloxy. U.S. Pat. No. 6,566,457 makes use of the same principal technology of hydrogenating a polymer in latex form in the presence of a ruthenium and/or rhodium containing catalyst in order to prepare graft polymers.\nJP 2001-288212 describes a further process for hydrogenating diene-based polymer latices. Latices of 2-chloro-1,3-butadiene (co)polymers are mixed with solutions or dispersions of catalysts in organic solvents which dissolve or swell the (co)polymers, and then contacted with hydrogen. The catalysts used are the so-called Wilkinson-catalysts having the formula MeCla(P(C6H5)3)b wherein Me is a transition-metal, Cl is chlorine, b is an integer and equal to or bigger than 1 and a+b is an integer less than or equal to 6. In the Examples a latex of poly(2-chloro-1,3-butadiene) rubber having a Tg of −42° C. and an average number weight molecular weight Mn of 150,000 was added to a toluene solution containing RhCl(PPh3)3 and Ph3P, and hydrogenated at 100° C. and 5.0 MPa for 2 hours to give a hydrogenated polymer with Tg=−54° C. and Mn=120,000.\nIn the Journal of Applied Polymer Science, Vol. 65, No. 4, 1997, 667-675 two processes for the selective hydrogenation of the C═C double bonds in nitrile-butadiene rubber (“NBR”) emulsions are described, which are carried out in the presence of a number of RuCl2(PPh3)3 complex catalysts. One of the processes is carried out in a homogeneous system, in which an organic solvent, which can dissolve the NBR polymer and the catalyst and which is compatible with the emulsion, is used. The other process is carried out in a heterogeneous system, in which an organic solvent which is capable of dissolving the catalyst and swelling the polymer particles but not miscible with the aqueous emulsion phase, is used. Both processes can realize quantitative hydrogenation of the C═C double bonds with the help of a certain amount of organic solvent to dissolve or swell the polymers.\nU.S. Pat. No. 6,696,518 teaches a process for selective hydrogenation of nonaromatic C═C and C≡C bonds in polymers with hydrogen in the presence of at least one hydrogenation catalyst comprising ruthenium and/or rhodium and at least one nonionic phosphorus compound capable of forming a coordinative compound with the transition metal wherein the hydrogenation catalyst is incorporated into the aqueous dispersion of the polymer without adding a solvent. Ru and/or Ru complexes or Ru and/or Ru salts are suitable catalysts. Examples of preferred nonionic phosphorus compound are PR3 or R2P(O)xZ(O)yPR2 [R represents e.g. C1-10alkyl, C4-12cycloalkyl, C1-10alkoxy, aryl(oxy) and F; Z is a bivalent hydrocarbon residue; x, y=0.1]. For this special case, an acrylic acid-butadiene-styrene copolymer latex was firstly prepared by radical polymerization of a mixture of monomers also containing ruthenium(III) tris-2,4-pentanedionate, which means that the Ru-salt was dispersed into the monomer aqueous solution as the catalyst precursor before the polymerization. After having obtained the aqueous polymer dispersion, Bu3P was added to the latex. The system was stirred for 16 h at ambient temperature followed by hydrogenation at 150° C. and 280 bar. The catalyst was synthesized in-situ, therefore no organic solvent was used to transport the catalyst. The hydrogenation is carried out in aqueous dispersions instead of in an organic medium, although the in-situ synthesized catalyst is oil-soluble.\nIn J. Molecular Catalysis Vol. 123, no. 1, 1997, 15-20 it is reported on the hydrogenation of polybutadiene (PBD), as well as polymers having styrene-butadiene repeating units (SBR) or having nitrile-butadiene repeating units (NBR) which are present in emulsions. Such hydrogenation is catalyzed by water-soluble rhodium complexes like e.g. [RhCl(HEXNa)2]2 (HEXNa=Ph2P—(CH2)5—CO2Na) and RhCl(TPPMS)3 (TPPMS=monosulphonated-triphenylphosphine). The process, however, is carried out in the presence of some organic solvent. Under the conditions employed the catalyst is extracted into the organic phase during reaction. This is attributed to the phase transfer properties of the complex which is rendered by the amphiphilic HEXNa ligand.\nIn Rubber Chemistry and Technology Vol. 68, no. 2, 1995, 281-286 it is described to use a water-soluble analog of the Wilkinson catalyst, i.e. RhCl(TPPMS)3 (where TPPMS represents monosulphonated-triphenylphosphine), for hydrogenation of nitrile rubber latex without any organic solvent and in the absence of any co-catalyst. The hydrogenation occurs under mild conditions (1 atm of hydrogen pressure, 75° C.) without coagulation of the latex, and up to 60% hydrogenation can be achieved. It is recognized that the hydrogenation was accompanied by an increase in gel content of the latex.\nJP 2003-126698 A discloses a process for hydrogenating unsaturated polymer latices using a water soluble catalyst containing group VIII metals or their compounds and hydrophilic organic phosphine ligands in the absence of organic solvents. In the process as disclosed in the Examples, 0.133 mmol (as Rh) chloro(1.5-cyclooctadiene)rhodium dimer was stirred with 0.372 mmol of P(C6H4SO3Na)3 to generate an aqueous solution of the complex catalyst. One part of such catalyst solution was mixed with five parts of butadiene rubber latex without organic solvent for hydrogenation. However, the highest degree of hydrogenation is only about 56% which is unsatisfactory for a scale-up to larger production volumes.\nIn the Journal of Molecular Catalysis A, Chemical, Vol. 231, no. 1-2, 2005, 93-101 it is reported to perform an aqueous phase hydrogenation of polybutadiene-1,4-block-poly(ethylene oxide) (PB-b-PEO) by using water-soluble Rh/TPPTS complexes, wherein TPPTS means P(C6H4-m-SO3−). The hydrogenation can be successful only because the PB-b-PEO has water-soluble parts within its polymer chains. In such a hydrogenation system, mixed micelles are formed by mixing the amphiphilic PB-b-PEO with added cationic surfactant dodecyl trimethyl ammonium chloride (DTAC) and n-hexane. Hydrogenation conversion can go up to 100% after one hour catalyzed by Rh/TPPTS complexes ([Rh]=10 ppm or less in aqueous phase) generated in situ from RhCl33H2O and TPPTS under 80 to 100° C. and 20 bar of H2. Their recycling experiment showed that the catalytic activity of the anionic catalytic system, Rh/P(C6H4-m-SO3−)3, remained high in a consecutive run. The success of this hydrogenation system is mainly due to the fact that PB-b-PEO is an amphiphilic starting material. Therefore, the water-soluble catalyst works for systems which use amphiphilic polymer material.\nJP 2003-096131 A discloses a method for the latex hydrogenation of conjugated diene polymers in the presence of Pd compounds as catalysts under basic conditions. NaOH is dissolved in palladium nitrate solution and mixed with a Mg silicate slurry to give a mixture (pH 12), which is then filtered, washed and dried to give a catalyst. Then, 0.45 part of the catalyst is added to 60 parts of 15% acrylonitrile-butadiene copolymer solution for hydrogenation at 50° C. with 5 MPa of H2 for 6 hours to obtain 95% degree of hydrogenation. Gel formation information is not reported.\nIn US 2006/0211827 A1 a process for selectively hydrogenating nitrile-butadiene rubber latex without organic solvent is disclosed in which RhCl(PPh3)3 is used as water-insoluble catalyst and PPh3 is used as water-insoluble co-catalyst. This process can achieve high degrees of hydrogenation and does not show gel formation. However, this method requires long reaction times and high loading of the transition metal for synthesis of the catalyst.\nIn summary, several technical routes have been attempted to hydrogenate C═C double bonds in polymers being present in the form of a latex, which include using hydrazine or a derivative of hydrazine as a reducing agent together with an oxidant like oxygen, air or hydroperoxide, directly using oil-soluble catalysts accompanied with a certain amount of organic solvents, and using catalysts containing water-soluble ligands. The hydrazine relevant route has encountered a significant gel formation problem, especially when a high hydrogenation conversion is required, and gel formation is not desired for post processing operations. In all prior art references relevant to using oil-soluble catalysts, a certain amount of organic solvent is still required in order to achieve a reasonable hydrogenation rate. The route relevant to using water-soluble catalysts has also encountered significant difficulty in overcoming the crosslinking problem.\nThe present invention therefore had the object to provide a novel and improved process allowing the hydrogenation of a diene-based polymer present as a latex without gel formation, simultaneously with a high degree of hydrogenation within acceptable short reaction times and with a reduced amount of the transition metal being used for synthesis of the catalyst."}
{"text": "1. Field of the Invention\nThe present invention pertains to an apparatus for producing frozen confections without the use of molds, and more particularly concerns an improved transfer conveyor for use in such apparatus.\n2. Description of the Prior Art\nIn U.S. Pat. No. 3,985,223 issued to Peter W. Forcella et al on Oct. 12, 1976, a conveyor for transferring frozen confections is described as including a large number of articulated confection gripper mechanisms, hereinafter referred to as pickup carriages. The pickup carriages are interconnected to form an endless chain, and such endless chain of carriages is driven in timed relation to the freezer conveyor of the confection producing system. A pair of gripping tongs are detachably mounted to each carriage. The tongs are made in several configurations adapted to the different types of frozen confections to be produced, such as, for example, cones, stickless bar confections, and stick-type bar confections. More particularly, the tongs have either prongs at their outer ends adapted to cooperate to seize a stickless bar confection therebetween, or they have yolk type abutments at their outer ends for grasping cones, or they have interdigitating teeth for handling stick type confections. As shown in the Drawings in the Forcella et al patent, the prior art tongs which are adapted to handle stickless novelty bars of ice cream or the like have staple-like wire configurations which terminate in tips of reduced diameter adapted to penetrate into the sides of the bar confection.\nAs disclosed in the aforementioned Forcella et al patent, flanged mounting members are affixed to the inner ends of the gripping tongs, and two dowel pins are attached to the flange of each mounting member. The ends of the arms of the carriages to which the tongs are detachably mounted have two apertures formed therein for receiving the dowel pins that are affixed to the tongs. One of the two pins extends upwardly of each tong flange, and a spring clip is pivotally secured on such projecting pin. The spring clip has a C-shape which enables it to be rotated on the projecting pin to a position where it clamps the tong and the arm together and retains the tong on the arm. The spring clips of the opposed tong members are of opposite geometrical configurations, that is to say, they are not interchangeable from side to side on each carriage.\nThe primary shortcoming of the confection pickup conveyor shown in the Forcella et al patent is that it is relatively time consuming and difficult to replace the tongs of each carriage to enable the pickup conveyor to handle a different type of confection, such as, for example, to remove the tongs adapted to handle stick-type bar confections and then reinstall tongs adapted to handle stickless bar confections. More particularly, it has been found that the removable spring clips used to engage the tongs on the tong receiving arms of the carriages must be relatively stiff in order to rigidly fasten the tongs to the arms, and such stiffness makes it necessary for one to use a screwdriver to pry off such spring clips. Furthermore, since the spring clips are of opposite geometrical configurations, it is necessary that such clips be stored separately, because otherwise one might have some difficulty in finding the appropriate spring clip for the appropriate side of the carriage. Such problem in changing from one type of tong to another is particularly serious because a frozen confection transfer conveyor may have from 100 to 150 carriages, each of which has a pair of tongs which must be removed and replaced with a new pair of tongs."}
{"text": "The present invention is related to an establishment for viewing an image where spectators can see an image having a deep and high realistic scenery and the like from a compartment which is provided therein or is provided upon a movable object.\nA conventional image viewing establishment such as a theater, a planetarium, an institution for having an exhibition, or the like, projects or exhibits an image at its maximum. Also, in order to produce parallax, the establishment cannot hide an edge of the image and its circumference intentionally. In a movie theater and so on, though the image is partitioned by a black curtain, it is clear that the curtain is a circumference of the image only and is not the benefit which produces parallax. That is, since a close scenery component which produces natural parallax had been eliminated intentionally between a circumference of the image and the image, even though it may be a large image, the image projecting establishment could not give a highly realistic feeling and vividness to the image.\nMoreover, when the spectators see the usual scenery from a window, the scope of the scenery that can be seen differs greatly by a distance from a window and their actual position. It is a big factor that this parallax makes a scenery of a window point of view relate to an actual object. But, in a conventional image display establishment, even though a point of view was moved therein, the scope of the image that can be seen did not change at all and did not give depth to the image at all."}
{"text": "1. Field of the Invention\nThe invention relates to a connector.\n2. Description of the Related Art\nU.S. Pat. No. 7,448,908 discloses a waterproof connector with first and second housings that are fit together by operating a lever. The first housing has a frame with concave accommodation parts and sub-connectors having terminal fittings are accommodated in the concave accommodation parts. A collective rubber stopper is accommodated in the concave accommodation parts to seal gaps between electric wires pulled out of the sub-connectors and inner peripheries of the concave accommodation parts. An engaging part is formed on the outer surface of the frame and defines a supporting shaft for pivotably supporting the lever. The lever has a cam groove that engages a cam follower of the second housing. Pivoting the lever generates a cam action between the cam groove and the cam follower that urges the first and second housings together.\nFit-on resistance between the first and second housings creates a reaction force that acts on the engaging part when a pivotal operation force is imparted to the lever. Thus, there is a fear that the outer wall of the frame will deform. Such a deformation will generate a gap between inner surfaces of the concave accommodation portions and the outer surface of the collective rubber stopper and could adversely affect the waterproof performance.\nThe invention has been completed based on the above-described situation. It is an object of the invention to prevent deterioration of the waterproof performance of a connector."}
{"text": "1. Field of the Invention\nThe present invention relates to a device for enhancing the collection efficiency of the Solid Phase Microextraction (SPME) method of sample concentration, collection and injection for analysis by gas chromatography or liquid chromatography. More specifically, it relates to a device for enhancing the collection efficiency of the SPME method, which enhances the flow of the analytes onto the sampling fiber.\n2. Description of the Related Art\nHistorically, the chief method of analyzing trace chemicals, also known as “analytes”, was to employ a syringe to collect a sample of the analyte and then inject it into a gas chromatograph (GC) or liquid chromatograph (LC). This resulted in the capture and injection of only small quantities of the analytes and thus yielded poor sensitivity. It was discovered that if large quantities of vapor or liquid analyte were drawn through a treated filter the components of interest would be concentrated on the filter paper. Solvents were used to selectively remove the chemicals of interest from the filter and a small portion of the solvent containing the traces was then injected into the GC or LC. Concentration of the sample via a filter was an improvement in sensitivity over straight analyte injection, but it added many processing steps to the analysis, which could increase errors.\nIt was further discovered that gas samples could be drawn through tubing containing packing material, such as Tenax, to concentrate the sample on the packing. After sampling, the tubing was connected to the sampling inlet of the GC. If the packing was rapidly heated it would drive off the trace chemicals into the GC for analysis. This process was simpler than the filter process and yielded similar sample concentration enhancements as the filter paper method without the solvent extraction problems. However, this process was only compatible with GC techniques and not LC techniques.\nCertain materials selectively adsorb and absorb materials based upon their chemical properties. For instance, silica gel will absorb water and then release it when heat is applied. Other extraction materials will release to certain solvents. Solid phase extraction (SPE) is a process which may employ selective adsorbents and solvents to concentrate the components of interest within the sample and selectively remove components that may interfere with later analysis. SPE can be a known effective alternative to liquid-to-liquid extraction in the analysis of aqueous samples. SPE reduces the consumption of high purity solvents and also reduces the time required to isolate the analytes of interest. However, the fact that SPE continues to use solvents has a number of disadvantages, notable among them being the need to take extensive precautions to guard against interference of the solvent with the analytes in the analyses.\nSolid phase microextraction (SPME) is a similar process which allows the concentration of volatile or nonvolatile compounds from liquid samples or from headspace gas without the complicated apparatus and solvents of SPE. SPME may use a fiber mounted within a hollow needle sheath of a syringe. The fiber acts as a sponge which captures and concentrates the analytes. The fiber can then be inserted into the heated injector of a GC where the analytes may be thermally desorbed from the fiber and made available for analysis by the GC. SPME can achieve detection limits down to the parts-per-trillion range for a wide variety of compounds.\nFIG. 1 is a cross-sectional view of one SPME device 10 in a retracted position and adjacent to a sample vessel 20. A thin fiber 30 may be coated with a substance that can enhance the selective absorption of the analytes, typically based upon the polarity of the molecules. Fiber 30 is attached to one end of a rod 40, with the other end of rod 40 being attached to plunger 50. Fiber 30 may be positioned inside of the needle sheath 60. Needle sheath 60 may be attached to cylinder 70. In the retracted position, cylinder 70 contains plunger 50 and partially contains rod 40, allowing them to potentially move longitudinally therein as a force is applied to plunger 50 along its main axis. SPME holder 80 is generally cylindrical in shape and includes a septum 90 at one end and a rod support and seal 91 at the other end. Needle sheath 60 containing fiber 30 is positioned inside SPME holder 80 and adjacent to septum 90. Sample vessel 20 may contain the material to be sampled, solid, liquid or gas, and may include a septum 100.\nIn operation, septum 90 of SPME holder 80 may be placed adjacent to septum 100 of sample vessel 20. Cylinder 70 can be depressed, thus injecting needle sheath 60 into sample vessel 20, as shown in FIG. 2. Next, plunger 50 is depressed which may extend fiber 30 through the free end of needle sheath 60, and thus may expose fiber 30 to the liquid or headspace gas within the vessel, as shown in FIG. 3. In this manner, fiber 30 is likely not damaged by unwanted contact. The analytes diffuse onto the surface of fiber 30 where they can be absorbed. Once enough time has elapsed to allow sufficient sample to be captured by fiber 30, fiber 30 is retracted into needle sheath 60 by retracting plunger 50 and needle sheath 60 is retracted via cylinder 70. Having the fiber containing the analyte sealed behind septum 90 prevents contamination of the fiber and may prevent escape of the analyte. Using a similar process, fiber 30 can be then be inserted into the GC sample inlet (not shown). The analytes are then removed from fiber 30 by thermal desorption and injected into a GC or LC for analysis. Since fiber 30 can be protected by septum 90 while retracted inside of SPME holder 80, this process protects fiber 30 while transporting field samples to a laboratory for the most detailed analyses.\nWhile the SPME process accomplishes its intended purposes of concentrating and collecting materials of interest, it suffers from a number of drawbacks. First, the SPME process is dependent upon diffusion to transport the analytes from solid or liquid on the bottom of the vessel to the fiber for collection. Since the time for diffusion and absorption of the analytes depends on many factors (e.g. the analytes themselves, as well as the type and thickness of any coating on the fiber), this process can take hours and even days to accomplish. Such long sampling times have been problematic because adsorption is a bi-directional process and, therefore, with a long sampling time, the analyte being adsorbed also begins to desorb from the fiber resulting in a non-linear response and making quantitization difficult. Therefore, techniques have been developed to reduce the absorption time.\nFor example, it is known that stirring and/or forceful agitation or heat may be applied to the sample vessel to reduce the absorption time. Stirring can be accomplished through the placement of a magnetic bar within the analyte and the use of a standard magnetic stirrer. Another method of agitation is to apply ultrasonic vibrations to the vial. However, agitation techniques can lead to damage to sample vials, damage to the fiber and also to damage to mechanical and electrical parts. Moreover, stirring techniques, whether they are magnetic, ultrasound or other methods, often increase the risk of contamination. Also, heat applied, either directly or as a side-effect of forceful agitation, leads to a rise in the temperature of the sample. Since adsorption efficiency is temperature-dependent, this introduces an unwanted variable into the analysis. SPME can be very effective in allowing the identification of trace compounds due to its ability to selectively concentrate the analyte. However, there are so many variable affecting the capture efficiency of the SPME that absolute quantitative analysis can be very difficult.\nIn the prior art, improvements on the SPME device have been proposed. For example, U.S. Pat. No. 7,131,341 to Wareham and Persaud discloses an instrument into which an SPME fiber and needle sheath is inserted. The holder/case contains a motor to present the fiber to the sample and retract the fiber into the needle sheath. This patent indicates that once the sample is on the fiber, the SPME can be either exchanged for a fresh one or immediately analyzed in the Wareham hosing. This patent is focused upon on-the-spot analysis of odor samples to determine if mold is within a wall or fungus samples to determine tree rot. This application is distinctly different than that of taking the SPME sample back to a lab for a more detailed analysis. First, if the SPME is removed from its housing for later analysis, as it is no longer sealed, it will likely leach sample and gather contamination. Though this patent discusses a tagging method, this tagging method appears incompatible with standard SPME holders, and thus leaves unaddressed a desirable next step of injecting the SPME into the laboratory analyzer for high resolution analysis. Further, this approach does not allow traceability of the analysis results to the quantity of air processed.\nIn order to overcome these problems, what is needed is a device to enhance the efficiency of the SPME process, the device permitting more rapid transport and absorption of the analytes to the fiber for collection. Further, the device must allow injection of the SPME into a full-scale laboratory analyzer. These features thus address and solve problems associated with conventional SPME systems."}
{"text": "Integrated circuits often contain registers for storing register states. A typical integrated circuit may hold hundreds or thousands of register states during normal operation of the integrated circuit.\nDuring debugging operations (as an example), it would be very time consuming to have to exhaustively monitor all of the register states on the integrated circuit. As another example, during retiming operations, it would be very time consuming to have to check the performance impact of moving every single register on the integrated circuit. As yet another example, it would be very inefficient to have to transfer all register states from one integrated circuit to another during a live migration event (i.e., an event where running of an application is transferred from one computing resource to another).\nIt is within this context that the embodiments described herein arise."}
{"text": "By variation of the four substituents on the phosphonium cation along with the available anions, we can obtain an enormous number of possible salts. Among them, the salts that melt below the normal boiling point of water, which are known as “ionic liquids (ILs)”, have attracted extensive attention. In some applications where they are used as phase transfer catalysts in chemical synthesis, electrolytes in energy storage systems, including lithium ion batteries, solar cells, actuators, and supercapacitors, and additives in medicaments and water treatment, quaternary phosphonium based ILs show superior performance as compared to quaternary ammonium based ILs. However, quaternary phosphonium based ILs have been much less studied. The slow progress on quaternary phosphonium based ILs can be attributed to the difficulty in synthesizing their starting materials, for example phosphine derivatives, and further the process for preparation of them.\nGenerally, quaternary phosphonium salts are synthesized by reacting tertiary phosphines with alkyl halides, followed by an anion-exchange process if the anion other than halide ions (X−) is desired.R1R2R3P+R4X→[R1R2R3R4P]+X−  (1)[R1R2R3R4P]+X−+M+A−→[R1R2R3R4P]+A−+M+X−  (2)[R1R2R3R4P]+X−+H+A−→[R1R2R3R4P]+A−+H+X−  (3)\nAlternatively, it is done by the addition of a metal salt M+[A]− with precipitation of M+X− or by displacement of the original ion by a strong acid MAI with release of H+X−. It deserves attention that those salts obtained by anion-exchange process are possibly contaminated with a small amount of X− ions unless the exchange reaction is fully completed.\nAs a synthesis process for a quaternary phosphonium salt with the anion of SO42−, for example, a quaternary phosphonium chloride has to be prepared first as shown in equation (1), and further a reaction of the quaternary phosphonium chloride with sulfuric acid as shown in equation (3) is required. During the anion-exchange process, the reaction equilibrium is shifted to the right side by continuous removal of volatile hydrochloric acid. In another case of preparing a quaternary phosphonium salt with the anion of BF4−, a quaternary phosphonium chloride is prepared as shown in equation (1), and followed by a reaction of the quaternary phosphonium chloride with NaBF4 in organic solvent such as acetone as shown in equation (2). When a silver salt is used instead of the alkali metal salt, the anion exchange reaction proceeds fast. In despite of the high cost, however, it is still difficult to obtain the desired quaternary phosphonium salt of high purity wherein the content of halide ion is very low.\nOn the other hand, the variation in properties between the quaternary phosphonium salts, even those with a same cation but different anions, is dramatic. The salts with anions such as NO3−, CO32+, PF6−, BF4−, C2O42−, Al2Cl7−, CH3COO−, CF3SO3−, C4H9SO3−, CF3COO−, N(CF3SO2)2−, N(C2F5SO2)2−, N(C4F9SO2)2−, N[(CF3SO2)(C4F9SO2)]−, and C(CF3SO2)3−, etc., may exhibit lower melting point, higher electrical conductivity, lower viscosity and/or stronger hydrophobicity as compared to the quaternary phosphonium halides. This sort of variation in physical and/or chemical properties gives rise to wider application in green chemistry and as novel electrochemical materials in various electrochemical energy storage systems. It is of great significance, therefore, to develop a new process of preparing high pure quaternary phosphonium salts with different anions.\nU.S. Pat. No. 4,892,944 discloses a process of preparing a quaternary phosphonium salt by reacting a tertiary phosphine with a carbonic acid diester to form a corresponding quaternary phosphonium carbonate and further mixing it with an acid to perform decarboxylaiton, as shown in equations below.R1R2R3P+Me2CO3→[R1R2R3PMe]+MeCO3−  (4)[R1R2R3PMe]+MeCO3−+H+A−→[R1R2R3PMe]+A−+MeOH+CO2  (5)\nAccording to the process, numerous quaternary phosphonium salts having various anions can be efficiently produced. To some extent the contamination of the anions could be avoided because in the presence of acid the alkyl carbonate ion is readily transformed to alcohol and CO2, which can be removed by simple method. However, this methodology is only applicable when tertiary phosphines are used as substrates. More readily available phosphine (PH3) and other phophine derivatives, including primary phosphines (RPH2) and secondary phosphines (R1R2PH), can not be alkylated by carbonic acid diester to give quaternary phosphonium compounds as shown in equation (4)."}
{"text": "The present invention relates to a refrigerator unit and/or a freezer unit having one or more storage shelves for the reception of refrigerated goods or frozen goods, with the appliance having one or more extension arrangements which are connected to the storage shelves and by means of which the storage shelves can be pulled out of and pushed into the appliance.\nPull-out storage shelves of refrigerator units and/or freezer units have the substantial advantage with respect to fixedly mounted, non-movable shelves that the refrigerated goods or frozen goods in the rear region of the shelves are comparatively easily accessible even when refrigerated goods and/or frozen goods are located in the front region of the shelves. A clearing away of the front region is thus not absolutely necessary. A disadvantage of previously known storage shelves consists of the fact, however, that the position of the shelves is fixedly set so that the appliances are comparatively inflexible with respect to the adaptability of the inner space to changing needs."}
{"text": "Gas generating processes based on the decomposition or burning of chemical propellants are frequently used for a number of purposes, e.g. for the inflation of inflatable devices, such as airbags, life boats or life vests, for operating pneumatically driven devices, or for use in fire-extinguish devices.\nKnown chemical methods for obtaining relative cool gases are generally based on the decomposition or the burning of solid materials in special units. These materials are generally shaped in the form of a solid block of material or as (loosely) packed powder, loose granules or loose tablets. The hot gases generated from the decomposition of these materials are in general cooled with the aid of special chemical cooling agents or by specific designed features such as heat exchangers. The high temperature burning gases are passed through the layer of the cooling agent or the heat exchanger and the temperature of the gases decreases as a result of the endothermal decomposition process of, or heat absorption by the cooling agent. Such processes are described for instance in GB-A 1,371,506.\nIn RU 2108282 it is observed that one of the drawbacks of the hitherto known art as cited above is the relatively complicated structure of these units. Another drawback is that the known gas generators did not allow or provide for the gases to be cooled below 150° C., which limits the applicability of such gas generators to systems that can withstand such high temperatures. Further disadvantages include the formation of relatively large quantities of undesired side-products, e.g. carbon monoxide (CO) or nitrogen oxides, a large mass and large size of the generator. For nitrogen gas generators RU 2108282 has the drawback of a reactive slag remaining after usage. This reactive slag requires controlled dismantling of the used gas generator, which make them less suitable for consumer products.\nWO 01/23327 proposes a gas generating device for effectively generating nitrogen gas of a low temperature, in order to overcome the above drawbacks.\nAccordingly, the gas generator of WO 01/23327 comprises at least one first body, comprising means for the generation of (nitrogen) gas, and at least one second body, comprising means for the generation of a neutralisation agent, wherein means are present for contacting the said neutralisation agent with the said first body, to neutralise the reactive reaction products (slag) from the generation of gas in the said first body, and wherein means are present for operating the generation of a neutralisation agent in the second body at a temporal and/or spatial interval from the generation of gas in the first body.\nWO 03/009899 relates to a cool oxygen chemical gas generator, which allows the production of oxygen gas at a temperature below 50° C. Although this generator may serve as an alternative to nitrogen in some applications, the reactivity of oxygen can be disadvantageous or make the generated gas unsuitable for a specific purpose. In general, the gas cannot be used in fire extinguisher applications, and may contribute to oxidation reactions, e.g. corrosion of materials that come in contact with oxygen.\nThere is a need for an alternative to a gas generator for nitrogen or oxygen cool gas.\nIt is an object of the invention to provide such an alternative, in particular such alternative that overcomes one or more drawbacks of the gas generators mentioned in the above identified documents.\nIt is in particular an object to provide a gas generator that is suitable to generate relatively cool gas, comprising no or a relatively low amount of toxic or hazardous components.\nIt has now been found possible to provide a gas generator suitable to generate relatively cool gas mainly consisting of a different gaseous molecule than nitrogen."}
{"text": "A reverse proxy server may be used to represent a secure content server to an outside client, preventing direct, unmonitored access to an intranet host server's data from outside the intranet. The reverse proxy server may also be used for replication; that is, multiple proxy servers may be associated with a heavily used host server for load balancing. In a system employing the reverse proxy server, an outside client trying to access a host server is directed to the reverse proxy server. Real content may reside on the host server, safely on one side the firewall. The reverse proxy server may reside on the other side the firewall, and appear to the client to be the host server.\nA typical reverse proxy server for the Internet may be configured such that the incoming Hyper Text Transfer Protocol (HTTP) requests intended for the host server are directed to the reverse proxy server. Upon receiving the HTTP request, the reverse proxy server for the Internet may determine whether it has a copy of the requested object stored locally. If so, the reverse proxy server for the Internet may respond to the HTTP request. If not, the reverse proxy server for the Internet may employ a Transmission Control Protocol (TCP) connection over the Internet to the intended host server. The reverse proxy server may then employ the TCP connection to send the HTTP request to the host server. If the request is successful, the host server may send the requested object to the reverse proxy server.\nUpon receiving the requested object from the host server, the reverse proxy server may perform content processing, also called content patching. Content patching commonly refers to the process of modifying links in the requested object, which may be a web page or other form of data. When the links are modified, the reverse proxy server may send the requested object to the client.\nThe above described process may work well for objects including exclusively static Hyper Text Markup Language (HTML) data. However, many web pages and other forms of HTTP data include dynamic links created by eXtensible Markup Language (XML)java script or java applets. Furthermore, minor problems with links such as syntax errors, missing symbols are generally not recognized by content patching applications resulting in incorrect or invalid links in the processed web page. Thus, it is with respect to these considerations and others that the present invention has been made."}
{"text": "This invention relates to portable seats, particularly to those designed for use in stadiums while watching sporting events with provisions for holding food and drink.\nThe current standard stadium seat consists either of an aluminum bench attached to concrete or just concrete itself, without any back support. Such stadium seats are very uncomfortable to sit in while watching sporting events such as football, baseball and soccer, which often last for several hours. As a result, spectators usually leave the sporting events with sore buttocks and aching backs. In addition to the standard stadium seat having no back support, such seats also have no place to hold food or drink purchased at the event. Furthermore, there is usually little or no room to place the food or drink in a crowded stadium itself. Therefore, spectators must hold their food and drink or invite spillage.\nPrior inventions have attempted to solve the aforementioned problems but none have done it all, particularly in the same or similar fashion as the instant invention. For instance, U.S. Pat. No. 4,474,407 by Nazar shows a portable seat with provision for storage only. Thus, the Nazar invention lacks a tray and also, the means of attaching the seat to the stadium seat to prevent tipping or sliding. U.S. Pat. No. 2,865,433 by Warner, discloses a stadium seat with article carrier, but the article carrier is designed for storage, and cannot be used as a tray. Furthermore, the Warner patent has no back support for the spectator. Another patent, U.S. Pat. No. 3,120,404 by Bramming shows a stadium seat that has a fold down lunch box rather than a tray. Contrary to the instant invention, the lunch box in this case must always be out even when not in use and is not useable as a tray. Thus, none of the prior art as represented by these patents contain devices that solve all the problems as does the instant invention.\nThe instant invention resolves the aforementioned problems by providing a comfortable seat comprised of a seat bottom portion consisting of a cushion made of foam or other material covered by vinyl or cloth, a connecting back support with a foam cushion covered by vinyl, a slidably removable tray to hold drinks and food and a spring loaded retractable metal latch to hold the seat onto the stadium seat so that the occupant will not fall over backwards while using it."}
{"text": "1. Field of the Invention\nThis invention relates to high friction surfaces for use in abrasive applications and the preparation of such high friction surfaces. In particular, the present invention relates to abrasive tools prepared by use of a matrix of braze paste and abrasive particles which have been oriented in a preselected pattern.\n2. Description of the Prior Art\nU.S. Pat. No. 3,918,217, hereinafter referred to as Oliver, discloses an abrasive tool comprising a tool blank having small steel balls secured thereto which have been armed with magnetically oriented carbide particles. The carbide particles are secured together and to the balls by braze metal. An Oliver-type tool is typically prepared, in part, by first securing a single layer of small steel balls to the surface of a tool blank. A magnet is then secured to the tool blank below the balls. Next, magnetizable carbide particles are sprinkled onto the small steel balls. The magnetic flux concentrations produced by the flux path through the small steel balls cause the carbide particles to collect on the outermost portions of the balls to form conical structures. A braze paste consisting of a binder and a braze alloy is then applied to encapsulate all of the elements of the assembly. Finally, the entire assembly is subjected to heating in a brazing furnace which bonds the braze alloy, the magnetizable particles and the balls into a unified structure. When the assembly has cooled, it may be used as an abrasive tool.\nThus, Oliver discloses a technique for preparing an abrasive tool which requires placement of protrusions on the external surface of a tool blank. These protrusions are the necessary surface for the formation of the conical structures of particles. The protrusions may take the form of small steel balls (as described) or shapes formed by machining a profile in the external surface of the tool blank. In either case, the tool blank which is the foundation of the ultimate product must be covered with protrusions. The small steel balls are expensive and require substantial labor to apply. If the protrusions are produced by machining a profile in the external surface of the tool blank, substantial skill and specialized equipment must be used.\nFurther, to apply carbide particles to protrusions located on a cylindrical structure, the structure must be fixtured and indexed to various rotary positions while the magnetizable carbide particles are applied. Preparing a tool in this manner is very complicated and prone to quality irregularities due to the variety of magnetic and gravitation forces which may be acting on the particles located in different rotary positions about the structure.\nFurther, the carbide particles form structures which align with the magnetic flux emanating from the magnetized tool blank. In the case of a cylindrical structure, the cones align with magnetic flux emanating radially (or perpendicular to a tangent drawn to the surface) from the circumferential portion of the wheel. The conical structures formed thereby will be symmetrical in all respects, and when used as a cutting tool will offer a negative rake angle to the workpiece. The Oliver technique is capable of producing only these symmetrical conical structures. If more aggressive rake angles are desired, the Oliver technique is not appropriate.\nFinally, as with any cutting or abrading tool, heat produced during cutting and abrading is conducted from the point of contact between the tool and the workpiece into the structure of the tool. When using an Oliver-type tool, heat passes from the conical structures through small steel balls and into the tool blank. Since these balls are brazed to the tool blank, they offer little more than point contact and consequently present a substantial impediment to the flow of heat."}
{"text": "1. Field of the Invention\nThis invention relates to novel substituted pyrrolopyridine derivatives which selectively bind to Corticotropin-Releasing Factor (CRF) receptors. More specifically, it relates to tetrahydro-5H-pyrido[2,3-b]indol-4-amines, 9H-pyrido[2,3-b]indol-4-amines, and 1H-pyrrolo[2,3-b]pyridin-4-amines, and their use as antagonists of Corticotropin-Releasing Factor in the treatment of various disease states.\n2. Description of the Related Art\nCorticotropin-releasing factor (CRF) antagonists are mentioned in U.S. Pat. Nos. 4,605,642 and 5,063,245 referring to peptides and pyrazoline derivatives, respectively. The importance of CRF antagonists is described in the literature, for example, as discussed in U.S. Pat. No. 5,063,245, which is incorporated herein by reference in its entirety. CRF antagonists are considered effective in the treatment of a wide range of diseases including stress-related illnesses, such as stress-induced depression, anxiety, and headache. Other diseases considered treatable with CRF antagonists are discussed in U.S. Pat. No. 5,063,245 and Pharm. Rev., 43: 425-473 (1991).\nInternational application WO 9413676 A1 discloses pyrrolo[2,3-d]pyriridines as having corticotropin-Releasing Factor antagonist activity. J. Het. Chem. 9, 1077 (1972) describes the synthesis of 9-Phenyl-pyrrolo[3,2-d]pyrimidines."}
{"text": "1. Field of the Invention\nThe present invention relates to a substrate processing method and a substrate processing system.\n2. Description of the Related Art\nIn a photolithography process in a fabrication process of, for example, a semiconductor wafer (hereinafter referred to as a xe2x80x9cwaferxe2x80x9d), a resist solution is coated onto a surface of the wafer, and the wafer undergoes a resist coating treatment for forming a resist film, an exposure processing for exposing a pattern to the wafer, a developing treatment for developing the wafer after the exposure and the like, so that a predetermined circuit pattern is formed on the wafer.\nAs the resist coating treatment, a dominant method today is a spin-coating method which discharges the resist solution to the center of the rotated wafer to diffuse the resist solution over the surface of the wafer.\nHowever, since the wafer is rotated a high speed according to the spin-coating method, high volume of the resist solution scatters from the peripheral edge part of the wafer, and therefore a lot of the resist solution is wasted. Additionally, since the unit is contaminated by the scatter of the resist solution, there are disadvantages that the unit needs to be cleaned frequently and so on.\nHence, instead of the spin-coating method which rotates the wafer, a method of discharging the resist solution onto the wafer while relatively moving the wafer and a resist solution discharge nozzle so that a path of a resist solution discharge part becomes a rectangular wave shape with large amplitude, which results in that the wafer is coated evenly with the resist solution in a plurality of parallel lines, that is, a method of coating of the so-called one continuous writing stroke is conceivable. In this coating method of one continuous writing stroke, there is a possibility that the surface of the resist film after the coating swells along the coating path of the resist solution, and hence it is preferable to use the resist solution having low viscosity which spreads easily over the wafer after the coating, and to flatten the coating film after completion of the coating.\nBy the way, it is recently demanded to further narrow the so-called xe2x80x9cedge cutting widthxe2x80x9d, which is at the peripheral edge part of the substrate and is not commercialized, and hence a technology for narrowing the edge cutting width needs to be developed in the above-described coating method of the so-called one continuous writing stroke as well. In order to narrow the edge cutting width, it is necessary to increase the viscosity of the resist solution so that the resist solution in the periphery of the wafer does not swell.\nHowever, when the viscosity of the resist solution is increased in the coating method of one continuous writing stroke, the resist solution discharged onto the wafer becomes hard to be spread and to be smoothed in flattening processing after that, and hence film thickness is thickened along the coating path of the resist solution, and an in-plane uniformity of the resist film is decreased. Namely, in the coating method of the so-called one continuous writing stroke, lowering the viscosity of the resist solution in order to narrow the edge-cutting width and maintaining the in-plane uniformity of the resist solution are inconsistent to each other.\nThe present invention is made in view of the above-described points, and it is an object of the present invention to narrow an edge cutting width while using a coating solution having high viscosity and to maintain an in-plane uniformity of a coating film, even when coating processing of a substrate such as a wafer is performed in the manner of the so-called one continuous writing stroke.\nIn view of the above object, a processing method of the present invention comprises the steps of coating a coating solution on a surface of the substrate while relatively moving a coating solution discharge nozzle and the substrate and discharging the coating solution onto the substrate from the coating solution discharge nozzle, exposing the substrate to a solvent atmosphere of the coating solution after the step of coating, and reducing pressure inside a container in which the substrate is housed after the step of exposing.\nAccording to another aspect of the present invention, a processing method of the present invention, for processing a substrate, comprises the steps of coating a coating solution on a surface of the substrate while relatively moving a coating solution discharge nozzle and the substrate and discharging the coating solution onto the substrate from the coating solution discharge nozzle, applying pressure inside a container in which the substrate is housed after the step of coating, and reducing pressure inside the container in which the substrate is housed after the step of applying pressure.\nAccording to the present invention, for example, after the coating solution is coated in the manner of the so-called one continuous writing stroke, the substrate is exposed to the solvent atmosphere, whereby the solvent adheres to the surface of the coating solution and the viscosity of the surface of the coating solution can be decreased. Thereafter, by reducing the pressure inside the container in which the substrate is housed and forming airflow inside the container, the surface of the substrate whose viscosity is decreased is smoothed and flattened. At the same time, the solvent is evaporated and the substrate is dried by the airflow. Thereby, even when the coating solution having high viscosity is coated, the in-plane uniformity of the film thickness of the coating film can be maintained. Therefore, even when the coating solution is coated in the manner of the so-called one continuous writing stroke, it is possible to use the coating solution having high viscosity, and to narrow the edge cutting width.\nWhen the pressure is applied inside the container in which the substrate is housed, instead of exposing the substrate to the solvent atmosphere, it is possible to prevent volatilization of the coating solution on the surface of the substrate and to allow the substrate and the coating solution to get to know well each other, thereby improving the flattening during reduced-pressure drying after that.\nA processing system of the present invention can perform the above-described processing method efficiently.\nNamely, a processing system of the present invention comprises a coating unit for coating the substrate with a coating solution, a solvent atmosphere unit for exposing the substrate to a solvent atmosphere of the coating solution, and a reduced-pressure drying unit for subjecting the substrate to reduced-pressure drying, wherein the coating unit comprises a coating solution discharge nozzle for discharging the coating solution onto the substrate, and a moving mechanism for relatively moving the coating solution discharge nozzle and the substrate, wherein the solvent atmosphere unit comprises a chamber for holding the substrate in a prescribed atmosphere, and a solvent vapor supply mechanism for supplying a solvent vapor of the coating solution of a prescribed concentration into the chamber through a supply pipe, and wherein the reduced-pressure drying unit comprises a pressure reducing mechanism for reducing pressure inside a container in which the substrate is housed. Therefore, it is possible to narrow the edge cutting width while using the coating solution having high viscosity and to maintain the in-plane uniformity of the coating film formed on the substrate.\nAccording to another aspect of the present invention, a processing system of the present invention comprises a coating unit for coating the substrate with a coating solution, and a reduced-pressure drying unit for exposing the substrate to a solvent atmosphere of the coating solution and thereafter subjecting the substrate to reduced-pressure drying, wherein the coating unit comprises a coating solution discharge nozzle for discharging the coating solution onto the substrate, and a moving mechanism for relatively moving the coating solution discharge nozzle and the substrate, and wherein the reduced-pressure drying unit comprises a container for housing the substrate and containing the substrate airtightly, a solvent vapor supply mechanism for supplying a solvent vapor of the coating solution of a prescribed concentration into the container through a supply pipe, and a pressure reducing mechanism for reducing pressure inside the container.\nAs described above, the processing system of the substrate includes either the coating unit which can perform the coating method of the so-called one continuous writing stroke, the solvent atmosphere unit which can expose the substrate to the solvent atmosphere and the reduced-pressure drying unit which can subject the substrate to the reduced-pressure drying, or the aforesaid coating unit and the reduced-pressure drying unit which can expose the substrate to the solvent atmosphere and subject the substrate to the reduced-pressure drying, and therefore, it can perform the above-described processing method of the present invention preferably. Therefore, it is possible to narrow the xe2x80x9cedge cutting widthxe2x80x9d while using the coating solution having high viscosity and to maintain the in-plane uniformity of the coating film formed on the substrate. Moreover, according to the latter processing system, the processing of exposing the substrate to the solvent atmosphere and the processing of subjecting the substrate to the reduced-pressure drying can be preformed in the same reduced-pressure drying unit, and hence it is possible to save time of carrying the substrate. Furthermore, it is possible to mount larger number of processing units in the processing unit, and hence its processing capacity can be improved."}
{"text": "1. Field of the Invention\nThe invention relates to a closure for an opening in the housing or cover of a piston pump or the like, in particular a high-pressure fluid pump, and especially a closure for a valve housing, the closure having an insertable inner cover, a loadable and relaxable pressure-applying member for the said cover, a supporting cover which can be screwed into a threaded portion of a housing or cover part, and a loading device for the pressure-applying member.\n2. Description of the Prior Art\nIn a known closure of this kind (German Patent Specification No. 10 72 031 corresponding to British Patent Specification No. 775314), a piston is provided in the supporting cover to form a pressure-applying member, and this piston can be loaded by a hydraulic pressure medium at the end that faces away from the inner cover. After the required final pressure is reached, a pressure screw in the supporting cover is screwed against the upper end of the piston so that this can also be held pressed against the inner cover when the pressure medium is relaxed.\nBecause of the purposes for which they are intended, the openings that have to be closed off in housing parts or in the cover of a piston or plunger pump are often of considerable size. They must, for example, permit a valve to be changed. At the same time the closure must be capable of withstanding the operating pressure, which, in some instances, is an alternating load. In the case of closure comprising screwed-together parts, there arises the danger of the screwed connections loosening in the course of time during operations. If the connection is not retightened in time, the threads may become flattened and the destruction of individual components may even occur."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor memory device and a refresh method thereof.\n2. Description of the Related Art\nA dynamic random access memory (hereinafter, referred to as DRAM) is configured as one of semiconductor memory devices. A memory cell used in the DRAM is ordinarily composed of a charge storage capacitor and a MOSFET for controlling input/output to/from the capacitor. Although the information stored in the memory cell is represented by a form of charges stored in the capacitor, since the charges decrease as a time elapses, a refresh operation is necessary to update the stored information at every predetermined time (refresh cycle).\nTo effectively execute the refresh operation, the semiconductor memory device has a plurality of types of refresh operation modes. The refresh operation modes are, for example, a CBR (CAS before RAS refresh) mode for starting a refresh operation by inputting signals of CAS (column-address strobe) and RAS (row-address strobe), which are supplied from the outside, at a predetermined timing, a self-refresh mode for automatically executing a refresh operation to all the memory cells at a regular interval when the memory device is continuously placed in a waiting state during at least a predetermined period of time, and the like.\nIn addition to above mentioned refresh operations, a kind of refresh operation is also executed when an ordinary read/write operation is executed. That is, the contents, which are stored in all the memory cells connected to a word line selected to execute the read/write operation, are supplied to the bit lines connected to the memory cells, respectively, and, after the potentials of the bit lines are amplified by a sense amplifier, respectively, the stored contents are restored in the respective memory cells.\nA current consumed in these refresh operations depends on time intervals (refresh cycle) of the refresh operations, and the refresh cycle must be increased to reduce the current to be consumed. However, when the refresh cycle is increased, the charges stored in each memory cell are reduced and written information may be lost. Further, a reduction rate of the charges stored in the memory cell depend on the temperature of the memory cell, and the charges tend to decrease faster at a higher temperature.\nAccordingly, various countermeasures are employed to effectively execute the refresh operation at a low current consumption. For example, Japanese Unexamined Patent Application Publication No. 05-036274 discloses a technology of a self-refresh mode for executing a first refresh operation at a short refresh cycle and thereafter executing it at a long refresh cycle. Further, Japanese Unexamined Patent Application Publication No. 2002-373489 discloses a technology for setting a different refresh cycle in an automatic refresh operation and in a self-refresh operation in response to a command input to a memory device. Further, Japanese Unexamined Patent Application Publication No. 2002-343079 discloses a technology for measuring a temperature of a semiconductor chip and changing a refresh cycle according to the temperature of the chip.\nAs described above, a reduction rate of charges stored in a memory cell of a semiconductor memory device depends on temperature. Accordingly, to effectively execute a refresh operation at a low current consumption, it is preferable to execute a refresh operation at a short refresh cycle when a semiconductor device has a high temperature and at a long refresh cycle when it has a low temperature.\nHowever, neither Japanese Unexamined Patent Application Publication Nos. 05-036274 nor 2002-373489 discloses an effective refresh method in consideration of a temperature change of a semiconductor memory device. Therefore, the problem of the effective refresh method in consideration of the temperature change of the semiconductor memory device remains unsolved. Further, the technology disclosed in Japanese Unexamined Patent Application Publication No. 2002-343079 copes with the temperature change of the semiconductor memory device. However, a temperature measuring unit must be newly provided to measure a temperature of a semiconductor chip, from which a new problem arises in that the configuration of the semiconductor memory device is made complex. Thus, a further improvement is desired to execute the refresh operation at a low current consumption in consideration of a temperature change of a semiconductor device."}
{"text": "1. Field of the Invention\nThe present invention relates to a ball mill for grinding materials such as cement clinker and, more particularly, to a partition diaphragm for use in a ball mill, the partition diaphragm being formed inside a body of the mill and having multi-grinding chambers.\n2. Prior Arts\nGenerally, plural grinding chambers, i.e., multi-grinding chambers are provided in a ball mill for the purpose of improving productivity while reducing energy cost per material thereby improving grinding efficiency as a whole, and in which a series of grinding processes are performed comprising the steps of primarily grinding a bulky material using a ball of large diameter (serving as a grinding medium) in a primary grinding chamber; filtering the material ground to a certain grain size; feeding the filtered material to a secondary grinding chamber; and grinding the fed material finely.\nFor that purpose, the partition diaphragm placed between one grinding chamber and the other performs a filtering action through a screen plate which is stretched over the partition diaphragm and permits only materials ground to a cetain grain size to pass through. In this filtering, if only one screen plate is employed, it is essential that this screen plate is resistant to the impact of the ball in the grinding chamber and that the clearance of the slits provided in the screen plate be about 4 mm. From the viewpoint of casting technique, it is not always easy to secure such a small clearance, and moreover if trying to secure a certain thickness considering sufficient resistance to the impact, there arise further disadvantages such as deflection or blinding of the slits.\nSeveral partition diaphragm have been developed and proposed to solve the mentioned disadvantages as disclosed in Japanese Laid-Open Publication (unexamined) No. 54-127060, the contents of which were further improved and disclosed in Japanese Patent Publication (examined) No. 60-8867. In this improved proposal, a primary screen plate of a partition diaphragm facing primary grinding chamber is formed of a relatively thick member provided with slits each having 20 mm in clearance so as to be resistant to dynamic contact with a ball of large diameter and secure void area. On the other hand, a secondary screen plate disposed in parallel to the mentioned primary screen plate with a rib therebetween is provided with slits each having 6 mm in clearance. Thus, a connecting chamber for communicating with the entire circumferential direction is formed by the mentioned two screen plates. The secondary screen plate can be formed to be relatively thin because this screen plate is not put in contact with the ball. Furthermore, the secondary screen plate can be provided with a large number of punched small holes in view of ventilation.\nIn the actual operation of the ball mill, it is an essential requirement to balance the level of material in the primary grinding chamber with the level of material in the secondary grinding chamber. In case of the mentioned prior art, the partition diaphragm is formed into a double structure, in which the slit clearance of the primary screen plate is 20 mm and that of the secondary screen plate is 6 mm. As a result of such an arrangement, the problem pertinent to the partition diaphragm comprised of only one screen plate having a slit clearance of 6 mm in clearance of slits are certainly solved, but from the viewpoint of practical operation, it is not always optimum to establish the grain size of the material passing from the primary grinding chamber to the secondary grinding chamber to be smaller than 6 mm in diameter. To be more specific, when carefully observed, it is found that the amount of material fed from the primary grinding chamber to the secondary grinding chamber passing through the filtering slits is excessive over the grinding capacity of the secondary grinding chamber, and that a number of relatively large size (approximately 6 mm) grains are mixed and left in the material fed to the secondary grinding chamber and discharged therefrom as they are without further grinding, which results in the application of a further load to the secondary grinding chamber. Such a disadvantage can be obviously recognized by checking the lowest level of the material presented in a vertical section of the first and second grinding chambers as shown in FIG. 5. That is, referring to FIG. 5, the level of the moving material immediately before the partition diaphragm of the primary grinding chamber is largely lowered, thereby the grinding efficiency at such a portion is also largely lowered. This negative influence is carried over to the next secondary grinding chamber, and such unbalanced movement of the material will eventually bring about a decline in the grinding performance of the entire ball mill. Such an essential problem as mentioned above still remains unsolved and has not been overcome by the mentioned proposal of a mere double structure of a partition diaphragm.\nIt may be an obvious improvement to establish a slit clearance of the secondary screen plate to be smaller than 6 mm. It is, however, not always easy to manufacture by casting a screen plate the slit clearance of which is smaller than 4 mm. Under the present state of casting technique, disadvantages such as blinding, deflection occur easily and it will be an excessive demand beyond the state of the art in casting to establish and maintain a constant slit clearance smaller than 4 mm. It is certainly possible to provide by punching an iron plate with a large number of circular holes the diameter of which is smaller than 4 mm. But in the case of such punched holes, being quite different from the screen plate manufactured by casting, there remain sharp edges on the punched holes. And it may be sometimes the case that grains of material are caught by such sharp edges and sticked thereto resulting in blinding of the holes thereby causing the operation of the entire plant to stop. This kind of problem occurs easily when the hole size (i.e., hole diameter) is smaller, and the stuck material may bring about corrosion of the screen plate which, in turn, needs troublesome works such as disassembly or replacement of the partition diaphragm. Frequent occurrence of such operation stoppage has a grave influence not only on the grinding process but also on the production activities of the entire cement plant as a matter of course."}
{"text": "1. Field\nApparatuses and methods consistent with exemplary embodiments relate to a user terminal device and a displaying method thereof, more particularly, to a user terminal a user terminal device including a display panel in a square shape and a displaying method thereof.\n2. Description of the Related Art\nWith the development of electronic technology, various types of user terminal devices have been developed and distributed. The size of user terminal devices has been decreasing, but the functions of the user terminal devices have increased and become more diverse. Therefore, the demand for user terminal devices has been increasing.\nA user terminal device may provide various contents such as multimedia contents and application screens upon the request of a user. A user may select a desired function using a button or a touch screen provided on a user terminal device. A user terminal device may selectively execute a program according to a user interaction between a user and the device, and display the result thereof.\nAs more diverse functions are provided by a user terminal device, there are various needs for different methods of displaying contents and/or user interfaces. That is, as the method of displaying content has changed and the type and function of the contents have been increasing, interaction methods such as merely selecting a button or touching a screen may be insufficient to perform various functions of a user terminal device.\nThus, there is a need of user interaction technology which enables a user to use a user terminal device more conveniently."}
{"text": "1. Field of the Invention\nThe present invention, relates to a device designed to show the driver at any time whether the vehicle is operating economically or not.\n2. Description of the Prior Art\nAmong the known devices which function similar to the present invention and are commercially available are those known devices which measure the vacuum in the motor inlet manifold. But this measurement lacks precision and the readout is difficult to interpret.\nThose economizers that are attached to the carburetor are not adaptable to the new fuel injected motors. Thus, they are of interest for only one category of vehicles, i.e. carburated vehicles.\nMeasurement and indication devices which measure the consumption of carburant are known. These sample several variables such as the speed of the motor or of the vehicle, the gas butterfly valve opening and the selected gear ratio. But these are instruments which give a composite indication of the consumption calculated as an aggregate of the preceding information to be displayed on the dash board in a numerical or analog form. In other words these indicators inform the driver of the average instantaneous consumption, but do not give precise indications as to the time of and cause for overconsumption and the means necessary to remedy such situations.\nIn effect when a measurement instrument gives quantitative information more precise than a simple alarm signal, the driver is not able to see or can only poorly interpret the many instrument readings, as he is preoccupied with driving the vehicle."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for rolling chain link fence. Specifically, the apparatus of the present invention provides for the rolling of chain link fence with the links in a collapsed state to produce rolls having a tighter configuration. With the apparatus of the present invention rolls are produced from a predetermined length of fencing and with these rolls having a smaller diameter than rolls produced by prior art apparatus.\n2. Description of the Prior Art\nIn the prior art apparatus has been used to provide for the rolling of predetermined lengths of chain link fence into rolls for shipment. With the prior art apparatus, the links of the chains link fence are held in an open position by the rolling apparatus and the individual rolls of chain link fence are accordingly made with the links in the open position. The prior art rolling apparatus therefore produces rolls having a predetermined diameter in accordance with the length of the chain link fence. The prior art rolling apparatus did facilitate the rolling of chain link fence but the links are held in an open configuration so that the diameter of the roll is large relative to the weight of the roll. When shipping rolls of chain link fence, it has been determined that the volume of shipping space is filled, or the height of the pile of rolls becomes too high before the weight limitations are exceeded. Therefore, it would be desirable to provide for the rolling of chain link fence into a tighter roll configuration having a smaller diameter so that more rolls of chain link fence can be fitted into a given volume."}
{"text": "In the conventional liquid crystal display (LCD), light-shielding matrixes and polarizers are generally formed in two independent layers, as a result, not only the production cost is increased but also the display device is thickened.\nQuantum dot materials have been used to manufacture a color filter (CF) substrate in an LCD device. Quantum dot is a semiconductor nanostructure for binding a conduction band electron, a valence band hole and an exciton in three spatial directions. The confinement can arise from electrostatic potential (generated by external electrodes, doping, strain, or impurities), interface of two different semiconductor materials (e.g., in self-assembled quantum dots), surface of semiconductors (e.g., semiconductor nanocrystals), or a combination thereof. Quantum dot is a nanoparticle formed by II-VI or III-V group elements. The particle diameter of the quantum dot is generally ranged from 1 to 10 nm. As electrons and holes are confined by quantum dots, continuous energy band structures become discrete energy band structures with molecular properties which can emit fluorescent light after excitation. Based on quantum effect, the quantum dot can be widely applied in the fields such as solar cells, luminescent devices and optical biomarkers.\nIn the display field, a white light-emitting diode (LED) backlight and conventional color filters (CFs) are generally matched with each other to form a color display light source which has the problems of low utilization rate, narrow color gamut of display, etc. When the quantum dots are adopted as CF materials, the display color gamut can be effectively improved."}
{"text": "Conventionally, from the point of view of the effective use and reuse of paper resources, there are writing instruments that use inks that become colorless when heated to a high temperature such as above 60° C., etc. Further, there are ink jet printers that have a recording section using inks that become colorless when heated and a heating section that is placed before said recording section and that makes the ink colorless (see, for example, Japanese Patent Application Publication No. 2001-302954).\nIn a situation such as reading out by scanning the personal information written in questionnaires and converting it to an electronic form, and managing it by storing it in a data base, etc., it is required to appropriately dispose of the paper after converting to an electronic form thereby preventing the leakage of personal information. In such cases, if the entering is done using a writing instrument using an ink such as the above that becomes colorless upon heating, it is possible to erase the entered part by heating the sheet after it has been read out and the information is converted into an electronic form, and hence it is possible to prevent the leakage of personal information. In addition, if only the entered part can be erased, it is also possible to use the sheet again as a questionnaire form that has not been filled in.\nHowever, in a work flow in which, after reading by a scanner and converting the information to an electronic form, the sheet is erased by heating in a different apparatus, there is no guarantee that the erasing operation is definitely done after reading out and converting to an electronic form. In addition, if the time of leaving the document after reading and converting to an electronic form until the erasing operation is carried out becomes long, the possibility of information leakage becomes correspondingly higher.\nThe present invention was made with the aim of solving the above problem, and the purpose of the present invention is to provide an image reading apparatus that can, after reading out and converting to an electronic form a document containing a part that is written using a pigment that becomes colorless when heated, quickly erase the part that is written in such an ink."}
{"text": "Current long term evolution (LTE)/LTE-A frequency division duplex (FDD) networks may be unable to support half-duplex (HD)-FDD wireless transmit/receive unit (WTRU) if the HD-FDD WTRU simply follows regular FDD WTRUs. The LTE/LTE-A may support both FDD and time division duplex (TDD) together. TDD operation may be similar to half-duplex FDD, such that a WTRU may perform either uplink transmission or downlink reception at one time. Because a TDD uplink (UL)/downlink (DL) subframe configuration is reused for HD-FDD WTRU in FDD network and all TDD-specific functions are supported, the scheduling flexibility may still be restricted, as few subframe configurations were defined in TDD. Similar to FD-FDD, HD-FDD may use at least two separate frequency bands for UL and DL transmissions, but UL and DL transmissions may not be supported simultaneously. In other words, HD-FDD may suffer from spectral inefficiency due to incomplete frequency band usage. Meanwhile, HD-FDD operation may be scheduled in the same manner as a TDD system, and may operate asymmetric DL-UL or UL-DL data requirements in a manner similar to TDD networks. Furthermore, in either DL or UL transmission, there may be a need for corresponding acknowledgement/non-acknowledgement (A/N) feedback.\nIn FD-FDD, in order to minimize the control signaling overhead, UL hybrid automatic repeat request (HARQ) may be synchronized so that UL retransmission occurs 8 subframes after the previous transmission of the same HARQ process, which may not give the UL scheduling flexibility to the WTRU. In HD-FDD, DL and UL transmission may be scheduled by eNB to avoid the conflict transmission. However, for some situations, there may be no available or scheduled source for a scheduling request (SR) or for when the WTRU has a SR to be transmitted but may be preempted by DL data reception."}
{"text": "Disposal and reuse of wastewater from raw municipal sewage is problematic. Stringent wastewater treatment requirements have been promulgated to protect human health, particularly in those areas having limited water supply or dense populations. For example, Title 22 of the California Administrative Code establishes stringent water reuse criteria where human contact is likely to occur with treated wastewater.\nTypically, wastewater is disinfected by chlorination or ultraviolet irradiation where the treated wastewater is discharged to inland surface waters. Disinfection of this type typically achieves complete destruction of pathogenic bacteria and substantial deactivation of viruses, but does not provide complete virus destruction. Viruses have been detected in secondary effluents.\nTitle 22 of the California Administrative Code is directed to tertiary treatment requirements. Viral monitoring is not specified in Title 22 because viruses typically occur in low concentrations in treated wastewater. Viral monitoring is expensive. Viral assays require special expertise. Laboratory procedures usually are off line and time consuming. Analytical costs are high. Therefore, instead of imposing measurements of viral concentrations, Title 22 sets forth a tertiary treatment system that consists of chemical coagulation, sedimentation, filtration, and disinfection where the public may exposed to the treated wastewater, as in a recreational impoundment.\nUnder the provisions of Title 22, turbidity of the treated effluent normally cannot exceed an average operating value of 2 NTU after final filtration, and cannot exceed 5 turbidity units more than 5 percent of the time during any 24 hour period. Chlorination after this level of treatment typically insures effective virus destruction sufficient for the protection of public health. Direct filtration with chemical addition is allowed as an alternative to the complete treatment systems specified in Title 22 where it has been demonstrated that the results of the two treatment systems are comparable and meet the appropriate criteria.\nIt has been determined that disinfection rates typically correlate well with wastewater particle size distributions and that the ability to inactivate an individual wastewater particle is a function of the size of the particle. Direct tertiary filtration alone usually does not enhance the rate of disinfection unless the particle size distribution of the settled wastewater is modified. Tertiary filtration systems that operate to remove larger size particles should safely reduce the long contact times and high chlorine dosages typically employed in wastewater reclamation processes. Accordingly, granular filtration media is almost universally required as a part of wastewater reclamation. Granular filtration is somewhat time consuming and can be the limiting factor for a wastewater treatment system.\nMasuda et al. U.S. Pat. No. 5,248,415 discloses an upward flow filtration apparatus that is said to be useful as a tertiary filter for wastewater treatment systems and to operate at a relatively high flow rate. An embodiment of the subject matter described in the Masuda patent is represented in FIG. 24 generally at 28 and is labelled Prior Art.\nThe filtration media described in the Masuda patent comprise a plurality of crimpy fibrous lumps. The fibrous lumps are disposed in the upward flow filtration apparatus 28 between first and second perforated panels 36 and 38, respectively. The wastewater flows in an upward direction through the fibrous lumps and suspended matter is captured by the individual fibrous lumps.\nThe first perforated panel 36 is immovably mounted within the apparatus and the second perforated panel 38 is movably mounted within the apparatus and spaced below the first perforated panel. The lower movable perforated panel 38, or bottom plate, is raised to compress the fibrous lumps to eliminate air gaps and to form a dense filter layer. The wastewater passes upwardly through the movable bottom plate and the filter layer and exits the top immovable plate 36. Fine solid materials in the upward flow are said to progressively adhere to the filter layer from the lower portion to the upper portion thereof in that order. With progressive filtration, resistance to filtration is increased. The movable bottom plate is lowered from time to time and is said to define a cleaning chamber when the filtration performance is reduced and it becomes necessary to clean the fibrous lumps.\nHowever, the apparatus described in the Masuda et al. patent with the movable bottom plate has some difficulties associated with it. A ram or screw 40 for moving the bottom plate passes through the wastewater, the media, and the top stationary plate 36. The screw decreases the amount of room available for the media and potentially causes some channeling through the media in the region of the screw. The media are constructed of a loose fiber and can become entangled in the screw as it turns. A seal is used where the screw passes through the top plate, further complicating the operation of the device.\nWhen the bottom plate is moved to compress the media, the lower layers of media become compressed. The compressed filter media in the lower portion of the bed is then the first portion of the media to contact the wastewater since the filter is operated in an upflow mode. The filter clogs up fast because large particles and fines are both trapped by the compressed initial layers of the filter media. The entire unit is shut down, usually before the upper filter layers are fully loaded, and the filter media is washed before the next cycle is begun.\nThe filter media is washed by moving the bottom plate downwardly away from the media to define a cleaning chamber. However, the flow rate of the wash water makes it difficult to achieve separation between the media and to obtain efficient cleaning. The Masuda device has typically required washing of the filter media on a frequent basis at a full flow rate of wash water equivalent to the flow rate of wastewater. Thus the overall efficiency of the apparatus described in the Masuda patent is greatly reduced.\nIt would be desirable to develop a filtration system suitable for tertiary wastewater treatment that substantially reduced or eliminated at least some of the problems associated with the Masuda device and yet provided a high rate filtration system as a suitable alternative to granular filtration media."}
{"text": "Telephony has advanced dramatically with the advancement of technology. Telephone communication was once limited to an analog public switched telephone network (PSTN), where the PSTN has been traditionally formed of two types of telephone carriers, local and long distance. The local carriers established local networks for subscribers to communicate within local regions, and the long distance carriers created networks between the local networks to enable subscribers of different local carriers to communicate with one another.\nOver time, mobile telephone networks were developed to enable subscribers to use mobile telephones. At first, the wireless networks and handsets were analog. Technology for the wireless networks was developed to provide digital wireless communications, which provided a clearer signal than analog wireless communications.\nAbout the same time that the digital wireless networks were developing, the Internet was also developing. The International Standards Organization (ISO) developed an Open Systems Interconnection (OSI) basic reference model in 1977 that currently includes seven different layers. Each of the layers provides protocols for certain types of operations. More specifically, the seven layers include: physical layer (Layer 1), data link layer (Layer 2), network layer (Layer 3), transport layer (Layer 4), session layer (Layer 5), presentation layer (Layer 6), and application layer (Layer 7). Each entity interacts directly with the layer immediately beneath it and provides facilities for use by the layer above it. The protocols on each layer enable entities to communicate with other entities on the same layer. The Internet initially provided for simple digital data to be communicated between users. One of the early communication application included email. However, as communications standards and protocols developed, the Internet matured to include more advanced communication applications, including voice over Internet protocol (VoIP).\nFIG. 1 is an illustration of a legacy telecommunications network that includes class 4 and 5 switches 102a-102n (collectively 102) and 104a-104n (collectively 104), respectively, connected to a signaling system #7 (SS7) network 106 (indicated as dashed lines). Historically, the class 5 switches 104 were generally configured to communicate via in-band signaling verses the use of SS7 signaling and operate to form a local network of subscribers within the network to place telephone calls to one another. The class 4 switches 102 were developed for long distance connections between the class 5 switches 104 at end offices (not shown). The class 4 switches 102, which are monolithic, are generally formed of multiple components, including a port, port cross-connect matrix, switch messaging bus with external signaling units, and call processing unit, as understood in the art. Class 4 switches are circuit based and utilize time division multiplexing (TDM) and are capable of terminating higher high-speed communications, including T1, T3, OC-3, and other four-wire circuit connections. As understood in the art, TDM is a synchronous communications protocol.\nThe SS7 network 106, which includes signal transfer points (STPs) 108a-108n (collectively 108), service switching points (SSs) on the class 4 and class 5 switches, and service control points (SCPs—not shown). The SS7 network is connected to the class 4 and 5 switches for providing and maintaining inter-switch call services between the switches. The SS7 network is used to signal out-of-band call setup, as out-of-band signaling is more secure and faster than in-band signaling. The call state changes of the inter-switch trunks of the class 4 switch are communicated to the adjoining switches via the SS7 network via a connection to the STP. To manage and route calls, the STPs 108 are used as an Inter-Switch messaging network whereby two switches control the trunking between the switches via messaging over the SS7 network provided by the STP switches that act as the inter-switch message bus. A call state machine of the class 4 switch provides control for routing traffic within the cross-connect matrix of the monolithic switch. The call state machine also provides call control signaling information to other switches via the connections to the STPs. The call control signaling information is routed via the STPs to other switches for call setup and tear-down. The call control information routed by the STPs contains pertinent information about the call to allow the terminating switch to complete various calls.\nTelephony has benefited from the development of the OSI model in a vast number of ways. One way has been through separating the call controller into a distributed cross-connect on an asynchronous network, such as asynchronous transfer mode (ATM or a Internet Protocol (IP) network. FIG. 2 is an illustration of a conventional telephony network 200 that includes a packet network 202. In one embodiment, the packet network 202 is an ATM network. Media gateways (MGs) 204a-204n (collectively 204) are media translation or conversion devices that modify and convert protocols between disparate communication networks. The media gateways 204, which are in communication with class 5 switches 206a-206n (collectively 206) are located at the edge of the packet network 202. The media gateways 204 convert TDM packets or streams 208 into packets, frames, or cells (collectively referred to hereinafter as “packets”) 210 and vice versa.\nThe packet network 202 operates independently as a distributed virtual media gateway port cross-connect for voice calls primarily due to one or more call control managers (CCMs) 212 located on the packet network 202. The call control manager 212 is in communication with the media gateways 204 and operates to control the media gateways 204 and provide instructions on how to rotate the packets 210 via far-end address allocations. By separating the call controller from the class 4 switches, the packet network 202 becomes, in effect, the virtual cross-connect of the switching system. The packet network 202 enables packets 210 that include voice data, commonly known as bearer packets, to be tagged with a destination address 214a and origination address 214b for enabling content data 214c to be properly routed from the origination media gateway 204b to the destination media gateway 204a. The media gateways 204 use of the packet network 202 is controlled by the CCM 212 and may communicate the packets 210 over the packet network 202 via IP addresses and virtual circuit (VC) or virtual path (VP) between the media gateways 204 to appropriately route the packets to the correct destination network node through the packet network 202. CCM 212 receives call state processing information from the media gateways 204 and signaling points, and processes the call state changes by using look up tables (not shown). The CCM 212 thereafter communicates packet addressing and state changes to the media gateways 204 to process the call.\nEthernet protocol was developed to provide for a computer network that enables multiple computers to share a common external inter-communication bus. Ethernet is generally used to provide for local area networks (LANs). Ethernet operates by communicating frames of data. While Ethernet operates well within a local environment (e.g., within a building) because Ethernet assumes that there is an known capacity of bandwidth associated with the bus standards set forth in the IEEE 802.3 standard that defines Ethernet. Ethernet is a shared environment, where co-utilization creates transmission errors called collisions. These collisions are detected by Ethernet cards in computers and a random re-transmission timer is used to avoid the next collision. Ethernet poses special problems for use in communications systems given it lacks dedicated bandwidth and time slots. The shared nature of an Ethernet network creates additional complexities in that the amount of available bandwidth can vary when used with wireless technologies.\nCommunication protocols transmitted over packet networks, such as ATM or IP networks, may utilize TDM based transmission facilities, which are synchronous as compared to Ethernet transmission facilities, which are asynchronous. Synchronous transmission protocols utilize a common clock and channel schema so that each device on the network operates synchronously with a dedicated path. Two types of “connection” state knowledge are present in a dedicated system, such as a TDM. Each channel has a dedicated amount of bandwidth and an error rate that is calculated from a common clock to determine path errors. The two types of connection state awareness functionality are provided by the channel itself and the common clock and data within a TDM header. The common clock provides for a determination of (i) a communications data rate from one end-point to another end-point and (ii) the data quality. Additionally, the TDM protocol includes “far end state” data in a header of a TDM frame to indicate whether there is a connection at the far end, thereby providing an indication of continuity along the communications path. Specifically, in-band end-to-end alarming allows the cross-connect devices to receive indications of continuity problems with other end-points. The in-band alarming is also provided for connection quality, where Bit Error Rate (BER) allows each end-point to know the quality of the data being received. Furthermore, bandwidth is always in use meaning that packets are synchronous, which that the far end knows exactly how many packets are to be sent and received in a given time period (e.g., one second). Computation of utilization is easily made by using the known bandwidth and multiplying it by the “seizure” time or the amount of time in use.\nPacket-based communications sessions lack circuit based connection state awareness indicators and clocking functionality to provide a session controller the ability to know the path connectivity state to efficiently manage making call handling decisions with anything other than ample bandwidth to setup and use sessions. This lack of connection path state awareness with the communications protocols, such as Ethernet and Internet protocol (IP) technologies, result in “gaps” in terms of being able to react to decaying transmission path quality and be sensitive to shared use of bandwidth. Most IP call controller solutions are founded on enterprise applications, where a single entity owns the network and scale of the network is relatively small. IP and Ethernet protocols lack the in-band path signaling, quality and use metrics to allow for this scale, or the ability to perform enhanced call handling with paths outside the governance of the packet network. Because packet communications are asynchronous, there is no common clock, and, thus, there is no way to know how many packets were transmitted, which, in turn, removes the ability to characterize transmission quality of the entire path, the amount of bandwidth available, or the amount in use. Further, packet networks are “converged” meaning they have both real-time and non-real-time bandwidth use. Currently, there is no in-band mechanism for determining real-time and non-real-time bandwidth use; having such information would allow for handling calls. It is commonly understood that proper connection operational assumptions are made by call control engines when the SS7 signaling path is properly operating (e.g., provisioned bandwidth is available) between end-points within the SS7 network. These operational assumptions are problematic as Ethernet, IP, and other data networks become oversubscribed and cause the packet network to become congested and prevent throughput. In cases where an end-point, such as a WiFi telephone, is mobile and bandwidth changes with signal strength (e.g., a WiFi telephone losing bandwidth as an individual walks away from a connection point antenna), the connection operation assumptions also fail to provide graceful call handling.\nOne available technique in packet networks to prevent oversubscription of real-time media traffic is through the use of call admission control (CAC) or the IP equivalent known as Resource Reservation Protocol (RSVP). CAC is primarily used to prevent congestion in voice traffic and is applied in the call setup phase to ensure there is enough bandwidth for data flow by reserving resources. To reserve bandwidth through the entire packet network, a CAC requires that the CAC procedure be performed at each point along a virtual circuit between two media gateways on which a call is to be routed, and often in a bi-directional fashion. While CAC functionality exists, the use of such CAC functionality is almost never applied because of the amount of time needed by the CAC procedure during call set up. For example, currently, CAC typically cannot operate over 40 calls per second and typical call set-ups on media gateways or class 4 switches may be 200 calls per second or higher.\nOne technique used to monitor the performance of IP session performance (i.e., after a call session has been established) is the use of the real-time control protocol (RTCP) as defined in IETF RFC 3550. RTCP collects statistics on a media connection, including bytes sent, packets sent, lost packets, jitter, feedback, and round trip delay. Other information may be provided in the RTCP packet using profile specific extensions. RTCP, which operates on a per session basis, is used for quality of service (QoS) reporting after termination of a session. The statistics information may be used, for example, to improve the quality of service by limiting data flow or changing CODEC compression. Utilization of the real-time QoS statistics, however, is limited to the specific session associated with the RTCP stream.\nAn emerging standard that is being developed for Ethernet performance measures is 802.1AG. This standard operates by generating and communicating an 802.1AG packet or “heart beat” over an Ethernet network segment. The 802.1AG packets are communicated via a Layer 2 Ethernet Virtual Circuit, such as a VLAN or Ethernet tunnel. At the ends and mid-points in Ethernet tunnels, 802.1AG packets are transmitted periodically over the Ethernet network to the far end. The Y.1731 protocol is utilized to calculate the number of data frames communicated between the 802.1AG packets. This configuration enables a performance measures (PM) to compute certain information about the performance of the path between the end-points on an Ethernet network. This combination of 802.1AG and Y.1731 enables the end points to be knowledgeable about the Frame Loss Rate (FLR), packet delay, and jitter in the path. This configuration is helpful to assist in monitoring performance of an Ethernet network path and diagnosing connectivity faults. However, the configuration falls short of providing the amount of real-time bandwidth in use or the total bandwidth in use. This information is useful to the proper management by a session controller handling calls during periods of flux in the packet transmission path, or the management of the real-time traffic.\nService providers often have trouble isolating and diagnosing network problems. To attempt to locate a packet loss problem along a node segment (i.e., a path between two network communications devices) over a network a probe that may be used to trace data packets being communicated over the node segment. This probe, however, is typically an external device from the network communications devices and operates to run a trace over an instant of time to determine network performance information, such as packet loss, jitter, and delay. An operator using the external probe may view results of a trace to diagnose the network communications problem. These results are not accessible to the network communications devices and cannot be accessed by network communications devices to alter network communications.\nTelecommunications switching systems today provide for Internet protocol (IP) communications between two end-points within a network or a different network to be terminated to a far end-point. Calls between two end-points are routed to the terminating end-point based on the address input at the originator. This address information is then relayed to a Call Control Manger (CCM) that screens, translates, and routes the call to the terminating subscriber or to another network to be terminated at a far end subscriber's end-point. The basic functionality of this process is widely known within the art and is used throughout telecommunications networks for voice calling.\nWithin the architecture of this switching system, calls to and from end-points are controlled by the CCM. The CCM may be located within a monolithic device in a TDM switch architecture or provided by an outboard computing device that controls the calls by using signaling that controls network based routing and switching devices located within the network. The latter device is known as soft-switch architecture.\nThe soft-switch architecture within an IP network controls call processing through use of signaling to and from the end devices and media gateways. One example of a protocol used for this IP signaling is Session Initiation Protocol (SIP). This protocol is currently used mainly with IP telephony, such as VoIP, and can be used as an access protocol between the end-user and the CCM and/or between the CCM of one network and the CCM of another network.\nAnother protocol used mainly between the CCM and a media gateway is the ITU-T H.248 protocol, commonly known as Megaco. This protocol is a control protocol that allows the CCM to control ingress and egress from/to the media gateway as calls are set up using a media gateway. Within a packet network framework, IP communications between two end-points (both access end-points and media gateways) are controlled by the signaling of the end-point to/from the CCM. The CCM provides authentication, screening, translations and routing based on information that is stored in the CCM and from the state of the end-points that the CCM controls.\nWithin the soft-switch architecture, call control can only be accomplished based on information possessed by the CCM or the on/off state of the devices that has an association with the CCM. While this configuration is fine in a static environment, packet networks are in a state of change at all times since the network itself can carry different types of information besides voice calls. One skilled in the art knows that a packet network is a converged network that can carry voice, data, and video all in a single path, and routing of calls within a packet network is not static and can vary significantly from call to call.\nBecause of packet network content communications variables, calls may encounter congestion and loss of voice quality based on latency, jitter, and packet loss. These content communications variables can affect any portion of a call at anytime based on the network elements usage at the time of the call. Unlike a TDM system where dedicated channels and circuits are provided, the CCM only has control of it own end-points. Other end units may attempt calls, computers may send/receive data without talking to the CCM, and other devices may require bandwidth while the original call is progressing, thus causing voice quality problems for the participants. In addition to these basic gaps, many physical layer 1 systems that are poor in regulating bandwidth, are being used for transmission facilities. WiFi, EVDO, 4G (Wax), DSL, and cable systems are all physical layer 1 technologies that demonstrate different bandwidth rates and management of their ability to modify available bandwidth as the Signal-to-Noise (SNR) ratios change.\nConventional soft-switching has not been designed to provide relief for callers when congestion, jitter or delay problems, such as those described above, are encountered. Since conventional CCMs can only determine call success based on connectivity to and from the calling parties, voice quality between two parties is not taken into consideration for call success.\nCommunication problems of in-band signals over packet networks are difficult to isolate. Currently, if a communication problem exists over a transmission path, there are few techniques to isolate the problem. One technique includes using an external probe to capture and decode packets, commonly know as a trace, traversing over a communications path to help isolate the problem. However, technicians generally only run the trace in response to a customer notifying a communications carrier of a communications problem. If a problem exists across packet networks operated by different carriers, a typical response by a carrier is to contact the other to determine if the other carrier can locate a problem in its network. In other words, locating an in-band communications problem over one or more packet networks is difficult as troubleshooting tools for such problems are limited to out-of-band performance metrics (PM) and are not available as in-band information via control or signaling paths.\nA problem that exists with current implementations of telephony over packet networks is that a call control manager does not have information about the bearer path. Traditionally, there was a linkage between transmission path state and the monolithic switch that essentially owned one end of that path where the in-band signaling and line characteristics were available and was an integral part of the information used by the CCM for call processing. As demonstrated in current implementations of VoIP, without knowledge of the bearer path, the call control manager may establish calls that result in poor voice quality or call setup failure.\nIn addition, IP Service gateways, such as a broadband remote access server (BRAS), functions to limit, commonly known as traffic rate shaping, each customer's DSL traffic to their purchased speed. There is no end-to-end signaling, outside of the embedded TCP flow control mechanism, used to adjust the bursting to eliminate packet loss. Rate shaping is a statically forced bandwidth constraint that alters the nature of a transmission path in the packet networks. This shaping coupled with commonly shared or “over-subscribed” bandwidth normally associated with trunking facilities between networks results in unknown transmission path states between media gateways servicing VoIP and other real-time services, such as Video on Demand (VOD).\nTraffic Quality of Service (QoS) management of packets is performed, where multiple flows aggregate into a smaller flows or channels. The application of Internet Protocol QoS is performed at the egress point where traffic is transmitted over a single link. Current traffic engines use the following information to make QOS traffic decisions. The decisions are assigning a Class of Service (CoS) and then acting upon that service to shape, restrict, or pass traffic to an egress point. The variables used to assign priority to traffic flows can be based on: entrance port (assign a whole port a CoS), virtual circuit in a port (assign a CoS to an Ethernet Virtual Circuit, LSP, etc.), priority bit marking of each packet (P bit), protocol type (assigning a CoS to specific types of packets or traffic), IP address and port (assigning a CoS to a whole IP address, or its port addresses), session identification (a HTTP, UDP, or other session addressed call), or otherwise. This priority marking information is used by service points, and shared links to implement QoS for the shared traffic flows. QoS and CoS types of information are made available at the point of aggregation where traffic management or QoS functions occur. However, the number of packets transmitted or lost in the packet stream elsewhere in the network is currently unavailable without the use of a session or path based protocol. These packet loss functions are generally not tracked by QoS mechanisms.\nIn current traffic rate shaping designs, the Internet may burst a packet stream to a DSL user when the packet network or Digital Subscriber Line Access Multiplexer (DSLAM) itself may not have sufficient bandwidth to accommodate the packet session. In a TCP-based session, the transmission rate is throttled down after packet loss is detected in the session. In VoIP, the packet loss is not counted by the use of Real Time Control Protocol (RTCP) signaling, but it is captured as the call progresses by the end points. RTCP, however, only considers performance of its own sessions and not the transmission path performance as a whole. In both cases, packets are sent over the packet network that get dropped in mid-path and will not make it to the customer premises equipment (CPE) and user. More importantly, there is no cross-session information about the packet loss and no whole path information available in-band.\nAlso, packet loss can be due to available bandwidth transmission rate fall-off, such as when a WiFi user walks away from a WiFi Access Point (AP) and loses RF signal strength, signal-to-noise ratio increases, or congestion increases due to many users concurrently accessing the AP. All of these types of conditions in the transmission paths can have severe impacts upon the ability to accomplish call processing and call management.\nAn Internet Service Provider (ISP) may provide different Internet connectivity speeds or data transfer rates based upon their service plans. For example, a user may purchase 1.5 Mbits/sec data transfer rate for a predetermined amount, such as 10 Mbits/sec data transfer rate or higher. In general, the transmission path is between the shared (trunked) BRAS resource and the DSLAM that is supplying. The normal amount of bandwidth consumption in the download direction from the network to the user is high as compared to the upstream direction. However, there is no correlated throttling mechanism in the IP web-server linked to user's ISP service plan that can be used to shape the packet transmissions. So, all of the network traffic is shaped at the BRAS typically based on the users purchased data transfer rates. Depending upon network conditions, some of this traffic may not make it to the DSL user since the BRAS has no knowledge of the IP service path from itself to the customer.\nA problem occurs when the BRAS does realize congestion on a packet network where packets are being dropped due to insufficient bandwidth. Some packets could be dropped in the packet network or at an aggregation device somewhere in the packet network. Currently, there is little intelligence that recognizes the dropped packets in the packet network. In fact, packet networks are designed to discard traffic based on QoS markings. This problem is made worse because transporting packets that will ultimately be dropped adds to congesting the network. The packets consume bandwidth until dropped.\nTransmission Control Protocol (TCP) was designed to work in a best-effort, packet store-and-forward environment characterized by the possibility of packet loss, packet disordering, and packet duplication. Packet loss can occur, for example, by a congested network element discarding a packet. Here, a microprocessor or memory of a network element may not have adequate capacity to address all packets routing into and out of the element. Packet disordering can occur, for example, by routing changes occurring during a transmission. Here, packets of TCP connection may be being arbitrarily transmitted partially over a low bandwidth terrestrial path and as routing table changes occur partially over a high bandwidth satellite path. Packet duplication can occur, for example, when two directly-connected network elements use a reliable link protocol and the link goes down after the receiver correctly receives a packet but before the transmitter receives an acknowledgement for the packet.\nAn embedded capability within TCP protocol is the TCP sliding window technique. The sliding window was developed and deployed as a flow control mechanism used to minimize the inefficiencies of packet-by-packet transmission. The sending of data between TCP enabled end-devices on a connection is accomplished using the sliding window technique. TCP requires that all transmitted data be acknowledged by the receiving host. The sliding windows method is a process by which multiple packets of data can be affirmed with a single acknowledgement."}
{"text": "Tractors and skid steers are common utility vehicles in the farming, construction and landscaping industries. The primary benefit of these vehicles is their great versatility and the wide variety of implements that can be attached to and used with the vehicle. After an initial investment in a tractor or skid steer, owners of these vehicles can add increased functionality by purchasing a wide variety of specialized implements including loaders, backhoes, blades, post hole diggers, forks and bale spears, to name a few. Because of the ubiquity and great versatility of tractors and skid steers, owners of these vehicles continuously seek new and improved implements to perform specialized tasks.\nPinching fingers are a good example of an implement that is especially useful for performing a specialized task. Pinching fingers are essentially claws that can be opened and closed around material or an item to be transported by the utility vehicle. Because the pinching fingers can be closed around the material or item, bulky and irregularly shaped items can be picked up and securely held while transported by the utility vehicle.\nAlthough they have many uses, pinching fingers are especially useful for transporting logs and clearing brush. Pinching fingers can be clamped around a log so that the log cannot roll when it is picked up, transported and loaded, increasing the safety of working with large logs. Pinching fingers can also be easily clamped around irregularly shaped items such as fallen trees or heavy farm equipment, making it easier to lift and move such items than with conventional implements such as a bucket or forks.\nPinching fingers, however, are not easily adapted for use with conventional tractors or skid steers. Most tractors and all skid steers have arms extending from the front of the vehicle for the operation of a loader or other implements. These implements can be raised and lowered by hydraulic cylinders attached to each arm and also rolled forward and backward by a separate cylinder or cylinders. This range of movement (up, down, forward roll and backward roll) is sufficient for operating most common implements including loaders, blades, forks and spears. Thus, two sets of hydraulic controls, one for raising and lowering the implement and one for rolling the implement forward and backward, are standard on most tractors and skid steers. Pinching fingers, however, also requires means for opening and closing the pincer claws.\nPrior art implements for pinching and grappling a load have been adapted for use with conventional tractors and skid steers. These implements typically utilize a third set of hydraulic controls (or other power means) for opening and closing the pincer claws or grappler. This solution, however, increases the cost of the implement because hydraulic cylinders, couplers and lines must be incorporated into the design of the implement. In the case of a tractor or skid steer that does not already have the necessary controls and lines installed, an entire hydraulic control and power system must be installed on the tractor or skid steer to use the implement. This added cost and complexity has prevented many potential users of pinching fingers and grapplers from adopting this useful implements. Thus, those of skill in the art have sought an implement that can be easily attached to a conventional tractor or skid steer without the need for separate controls and power means for opening and closing the implement.\nIn addition to the problem of requiring a separate set of controls, prior art designs have not been easy to install and remove on a conventional tractor or skid steer. To maximize the versatility and thus the utility of these vehicles, implements must be designed so that they can be quickly interchanged with other implements. Thus, those of skill in the art have sought pinching fingers with simple mechanical means for installing and uninstalling the implement on a conventional tractor or skid steer.\nPrior art designs also typically suffer from another deficiency in their operation. To provide structural stability and strength, the top set of pincer claws on a conventional grappler are constructed as a unified or solid piece that clamps in opposition to and against a solid bottom set of pincer claws, forks or a bucket. When grasping an item of uniform cross section, this design is adequate for providing uniform pressure at each claw to secure the item. However, when a solid item to be grasped has a tapered or irregular cross section, one claw may be restricted from closing on the item when the other claw closes on a larger portion of the item, causing the item is held by only one claw and not the other. In this instance, the item cannot be safely transported without the danger of the item moving and shifting during transport. Thus, those of skill in the art have sought a grappler and/or pinching fingers with independent claws that can be closed on an item even if another claw is restricted from closing, but without sacrificing the structural stability or strength of the implement.\nAccordingly, an object of the present invention is to provide pinching fingers for a tractor, skid steer or the like that can be used by the vehicle without the need for an additional set of controls and power means for opening and closing the pinching fingers.\nA further object of the present invention is to provide pinching fingers that are easy to install and remove on a conventional loader or skid steer.\nYet another object of the present invention is to provide pinching fingers that can be quickly interchanged with other implements on a conventional loader or skid steer.\nAn additional object of the present invention is to provide pinching fingers that may move and grasp independently so that both claws securely grasp a solid item of tapered or irregular cross section.\nA still further object of the present invention is to provide pinching fingers attachment with sufficient structural stability and strength to grasp, lift and transport heavy items including logs, trees and other items.\nStill another object of the present invention is to provide an improved vehicle with a pinching fingers attachment.\nYet another object of the present invention is to provide an improved method for using a pinching fingers attachment with loaders, skid steers and the like.\nA further object of the invention is to provide a loader or skid steer with pinching fingers that are simple and easy to operate.\nFinally, an object of the present invention is to provide a pinching fingers attachment that is economical to manufacture, durable and refined in appearance."}
{"text": "The disclosure of Japanese Patent Application No. 2001-11842 filed on Jan. 19, 2001 including the specification, drawings and abstract is incorporated herein by reference in its entirety.\n1. Field of Invention\nThe invention relates to a reclining device. More specifically, the invention relates to the fitting of a seat back and a seat cushion.\n2. Description of Related Art\nA reclining device of a certain type is disclosed, for example, in Japanese Patent Application Laid-Open No. 8-253063. This reclining device is designed such that a cam body presses a plurality of slide pawls each having outer teeth on the front end side against an annular ratchet having inner teeth on the inner periphery side so that the slide pawls are brought into mesh with the ratchet and locked with respect thereto. Further, the operation of pressing the slide pawls against the ratchet is canceled by applying an arbitrary operational force to the cam body so that the slide pawls are retreated from the ratchet and brought out of mesh therewith.\nThe reclining device of this type is designed to prevent the ratchet from rotating relative to a holder when the slide pawls are locked with respect to the ratchet and to allow the ratchet to rotate relative to the holder when the slide pawls are unlocked from the ratchet. The reclining device of this type is disposed between a rear end portion of a seat cushion and a lower end portion of a seat back, thus constituting a vehicular reclining seat. For example, the reclining device of this type is designed such that the seat back is longitudinally tiltably supported with respect to the seat cushion by mounting the ratchet on the lower end side of the seat back and the holder on the rear end side of the seat cushion.\nIt is to be noted herein that the reclining device of this type confronts the following problem. If a great circumferential turning force is applied to the seat back constituting the reclining seat, the ratchet is also exposed to a great circumferential turning force, which acts on the slide pawls as radial and tangential components in the circumferential direction. The slide pawls are displaced both radially inwardly and circumferentially. As a result, the slide pawls may be brought out of mesh with the ratchet. Thus, the reclining device of this type requires that the slide pawls securely mesh with the ratchet when the slide pawls are locked with respect to the ratchet.\nHowever, the reclining device according to the related art as disclosed in the aforementioned patent publication is designed to bear each of the slide pawls only at a single point, namely, at the center of a rear end portion thereof. Hence, if a great circumferential turning force is applied to the slide pawls from the ratchet, each of the slide pawls sways circumferentially upon its bearing portion for the cam body and becomes quite unstable although it is still maintained in mesh with the ratchet. Then, the slide pawls may be retreated from the ratchet, disengaged therefrom, and unlocked therefrom.\nThe invention thus improves the reclining device of the above-mentioned type such that the slide pawls are stably maintained in mesh with the ratchet and cannot be disengaged therefrom easily even if a great circumferential turning force is applied to the slide pawls from the ratchet. To achieve this object, the reclining device of the invention adopts a construction capable of securely maintaining the slide pawls in mesh with the ratchet, taking into account the configuration of the slide pawls and the manner in which the cam body bears the slide pawls.\nA reclining device according to one exemplary aspect of the invention comprises a plurality of pawls, an annular ratchet, a cam body, and a holder. Each of the pawls has outer teeth formed on the front end side, first and second leg portions formed at opposed end portions, and a bearing portion formed between the first and second leg portions. The first and second leg portions extend toward the rear end side. The annular ratchet has inner teeth formed on the inner periphery side. The cam body comes into abutment at least on the bearing portion of each of the pawls and possibly on the leg portions and the bearing portion of each of the pawls. The cam body presses the pawls toward the ratchet so that the pawls are brought into mesh with the ratchet and locked with respect thereto and stops pressing the pawls so that the pawls are retreated from the ratchet and unlocked therefrom. The holder holds each of the pawls such that the pawl can move toward and away from the inner teeth of the ratchet, and rotates concentrically with respect to the ratchet. The cam body, which operates by an arbitrary operational force, releases the pressing force applied to the slide pawls and brings them out of mesh with the ratchet. The reclining device is designed to prevent the ratchet from rotating relative to the holder when the slide pawls are in mesh with the ratchet and to allow the ratchet to rotate relative to the holder when the slide pawls are out of mesh with the ratchet."}
{"text": "Field\nCertain aspects of the present disclosure generally relate to neural system engineering and, more particularly, to systems and methods for monitoring a neural network with a shadow network.\nBackground\nAn artificial neural network, which may comprise an interconnected group of artificial neurons (i.e., neuron models), is a computational device or represents a method to be performed by a computational device. Artificial neural networks may have corresponding structure and/or function in biological neural networks. In some cases, artificial neural networks may provide innovative and useful computational techniques for certain applications in which traditional computational techniques are cumbersome, impractical, or inadequate. Because artificial neural networks can infer a function from observations, such networks are particularly useful in applications where the complexity of the task or data makes the design of the function by conventional techniques burdensome. Still, like most systems, the neural networks may fail or become unstable. Thus, it is desirable to provide a neuromorphic receiver that includes a shadow neural network for monitoring the activity of a target neural network. The shadow neural network may improve the debugging and/or monitoring of the target neural network."}
{"text": "Hitherto, a heat-sensitive recording material which utilizes a coloring reaction of a color former and a color developer which forms a color upon contact with the color former and provides a color image by contacting both coloring materials by the action of heat is well known.\nSuch heat-sensitive recording materials are comparatively inexpensive and can be used on a recording equipment that is compact and requires fairly easy maintenance. Because of these advantages, heat-sensitive recording materials are extensively used as recording media in a wide range of applications including facsimile equipments and various types of computers. One major problem with heat-sensitive recording materials is, however, their low resistance to fingerprints or solvents; if the recording layer comes in contact with a sebum of human being or a solvent, the color density of characters to be attained is decreased or an unwanted coloration (\"background fog\") occurs.\nWith a view to solving this problem, several methods have been proposed; in one method, a heat-sensitive recording layer is coated with an aqueous emulsion of a resin having film-forming properties and resistance to chemicals (Japanese Patent Application (OPI) No. 128347/79); and in another method, a heat-sensitive recording layer is coated with a water-soluble high molecular weight compound such as polyvinyl alcohol (Japanese Utility Model Application (OPI) No. 125354/81). The term \"OPI\" as used herein means \"unexamined published application\". However, the methods so far proposed either suffer from additional problems or find themselves incapable of achieving the intended results to satisfactory levels. For instance, in the method in which an aqueous resin coating is applied to the heat-sensitive recording layer, the temperature for drying has to be limited to a certain level in order to avoid unwanted coloration of the recording layer due to high-temperature drying, whereby curing of the resin layer inevitably becomes insufficient for preventing its sticking to a recording head during the recording. In order to avoid this problem, a method in which a heat-sensitive recording layer is coated with a resin component capable of curing upon exposure to electron beams and the coated resin component is cured upon exposure to electron beams is proposed. However, the resulting heat-sensitive recording material is still unsatisfactory in terms of preservability of recorded images. In addition, there may also be problems that the electron beam-curable resin layer causes coloration of the heat-sensitive recording layer just after it has been coated or causes fading of recorded images.\nAttaching heat-sensitive recording materials to commercial goods has recently become a widespread practice, and in order to enhance the aesthetic appeal of the recording materials to thereby increase their commercial value as labels, it is desirable that the recording materials have either a high gloss or matted surface characteristics. This requirement is not fully met by the recording layer of existing heat-sensitive recording materials unless they are provided with an overcoat layer. However, the gloss attained by coating a heat-sensitive recording layer directly with an electron beam-curable resin layer as described above is not high enough to provide an increased commercial value for the product."}
{"text": "Deep brain stimulation (DBS) is a methodology wherein therapeutic electrodes are implanted within the brain to deliver timed impulses to desired nerve centers, and can be used to treat a variety of disorders, in particular movement disorders such as Parkinson's disease and dystonia. A major challenge with DBS relates to determination of where to place therapeutic electrodes: while medical imaging (e.g., magnetic resonance imaging) can provide a starting point for identification of placement locations, since the therapeutic electrodes must be precisely placed (e.g., in those regions of the brain giving rise to muscle tremors), it is necessary to obtain a more detailed map of brain structures. A common technique is to advance a needle-like probe (or multiple such probes) into the brain, with each probe bearing one or more reading microelectrodes. The reading microelectrodes measure neuronal activity, as indicated by extracellular action potentials, which are in essence voltage spikes caused by neuronal firings. As readings are taken, the patient—who is generally awake during the procedure—may be requested to perform some action (e.g., move an arm or leg), thereby provoking neuronal firings. By looking at the location of the probe (i.e., probe location and depth) and the characteristics of the measured neuronal spikes (e.g., shape, frequency, etc.), the regions of the brain traversed by the probe can be mapped: spike characteristics can be correlated with those known to exist in certain portions of the brain, changes in spike characteristics can indicate interfaces between different regions of the brain, and so forth. Additionally, stimulating input (voltage) pulses can simultaneously be delivered to regions of the brain which are candidates for electrode implantation to determine their physiological effect (e.g., whether tremors are reduced, whether the patient experiences some change in feeling, etc.), with such stimulating input pulses being delivered via an input point on the probe or via a separate electrode spaced from the probe. Once the map of the brain is generated and candidate locations for implantation of therapeutic electrodes are identified, therapeutic electrodes may be permanently implanted, with the therapeutic electrodes being connected to a power supply which delivers an input signal suitable to reduce or eliminate tremors, or to attain some other desired effect.\nHowever, the process of mapping the brain can be a difficult one. It can be extremely difficult to discern neuronal spikes from background noise (and from any stimulating input pulses) within these data, particularly owing to the wide variety of characteristics neuronal spikes may have; spikes can vary widely in their shape, amplitude, period, frequency, and so forth. The same neuron can even generate different spike readings over time, both owing to variability in the neuron itself and owing to factors such as nearby pulsing blood vessels creating small changes in neuron-to-probe spacing. Experienced neurologists and others can over time gain skill in identifying the neuronal spikes of individual neurons from probe readings, but since probes may contain large arrays of reading microelectrodes, thereby generating multiple streams of reading data, it is virtually impossible for human operators to successfully process all of the generated data. It would therefore be useful to have improved methods available for identifying neuronal spikes, in particular methods which require minimal human review and supervision, and which might therefore be suitable for use in expert systems and other automated or semi-automated systems for probe data review."}
{"text": "The goal of plant breeding is to combine in a single variety or hybrid various desirable traits. For field crops, these traits may include resistance to diseases and insects, tolerance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination and stand establishment, growth rate, maturity, and plant and ear height, is important.\nField crops are bred through techniques that take advantage of the plant's method of pollination. A plant is self-pollinated if pollen from one flower is transferred to the same or another flower of the same plant. A plant is cross-pollinated if the pollen comes from a flower on a different plant.\nPlants that have been self-pollinated and selected for type for many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny. A cross between two different homozygous lines produces a uniform population of hybrid plants that may be heterozygous for many gene loci. A cross of two plants each heterozygous at a number of gene loci will produce a population of hybrid plants that differ genetically and will not be uniform.\nMaize (zea mays L.), often referred to as corn in the United States, can be bred by both self-pollination and cross-pollination techniques. Maize has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in maize when wind blows pollen from the tassels to the silks that protrude from the tops of the ears.\nA reliable method of controlling male fertility in plants offers the opportunity for improved plant breeding. This is especially true for development of maize hybrids, which relies upon some sort of male sterility system. There are several options for controlling male fertility available to breeders, such as: manual or mechanical emasculation (or detasseling), cytoplasmic male sterility, genetic male sterility, gametocides and the like.\nHybrid maize seed is typically produced by a male sterility system incorporating manual or mechanical detasseling. Alternate strips of two maize inbreds are planted in a field, and the pollen-bearing tassels are removed from one of the inbreds (female). Providing that there is sufficient isolation from sources of foreign maize pollen, the ears of the detasseled inbred will be fertilized only from the other inbred (male), and the resulting seed is therefore hybrid and will form hybrid plants.\nThe laborious, and occasionally unreliable, detasseling process can be avoided by using cytoplasmic male-sterile (CMS) inbreds. Plants of a CMS inbred are male sterile as a result of factors resulting from the cytoplasmic, as opposed to the nuclear, genome. Thus, this characteristic is inherited exclusively through the female parent in maize plants, since only the female provides cytoplasm to the fertilized seed. CMS plants are fertilized with pollen from another inbred that is not male-sterile. Pollen from the second inbred may or may not contribute genes that make the hybrid plants male-fertile. Seed from detasseled fertile maize and CMS produced seed of the same hybrid can be blended to insure that adequate pollen loads are available for fertilization when the hybrid plants are grown.\nThere are several methods of conferring genetic male sterility available, such as multiple mutant genes at separate locations within the genome that confer male sterility, as disclosed in U.S. Pat. Nos. 4,654,465 and 4,727,219 to Brar et al. and chromosomal translocations as described by Patterson in U.S. Pat. Nos. 3,861,709 and 3,710,511. These and all patents referred to are incorporated by reference. In addition to these methods, Albertsen et al., of Pioneer Hi-Bred, U.S. Pat. No. 5,432,068, have developed a system of nuclear male sterility which includes: identifying a gene which is critical to male fertility; silencing this native gene which is critical to male fertility; removing the native promoter from the essential male fertility gene and replacing it with an inducible promoter; inserting this genetically engineered gene back into the plant; and thus creating a plant that is male sterile because the inducible promoter is not \"on\" resulting in the male fertility gene not being transcribed. Fertility is restored by inducing, or turning \"on\", the promoter, which in turn allows the gene that confers male fertility to be transcribed.\nThere are many other methods of conferring genetic male sterility in the art, each with its own benefits and drawbacks. These methods use a variety of approaches such as delivering into the plant a gene encoding a cytotoxic substance associated with a male tissue specific promoter or an antisense system in which a gene critical to fertility is identified and an antisense to that gene is inserted in the plant (see: Fabinjanski, et al. EPO 89/3010153.8 publication no. 329,308 and PCT application PCT/CA90/00037 published as WO 90/08828).\nAnother system useful in controlling male sterility makes use of gametocides. Gametocides are not a genetic system, but rather a topical application of chemicals. These chemicals affect cells that are critical to male fertility. The application of these chemicals affects fertility in the plants only for the growing season in which the gametocide is applied (see Carlson, Glenn R., U.S. Pat. No. 4,936,904). Application of the gametocide, timing of the application and genotype specificity often limit the usefulness of the approach.\nDevelopment of Maize Inbred Lines\nThe use of male sterile inbreds is but one factor in the production of maize hybrids. Plant breeding techniques known in the art and used in a maize plant breeding program include, but are not limited to, recurrent selection, backcrossing, pedigree breeding, restriction length polymorphism enhanced selection, genetic marker enhanced selection and transformation. The development of maize hybrids in a maize plant breeding program requires, in general, the development of homozygous inbred lines, the crossing of these lines, and the evaluation of the crosses. Pedigree breeding and recurrent selection breeding methods are used to develop inbred lines from breeding populations. Maize plant breeding programs combine the genetic backgrounds from two or more inbred lines or various other germplasm sources into breeding pools from which new inbred lines are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred lines and the hybrids from these crosses are evaluated to determine which of those have commercial potential. Plant breeding and hybrid development, as practiced in a maize plant breeding program, are expensive and time consuming processes.\nPedigree breeding starts with the crossing of two genotypes, each of which may have one or more desirable characteristics that is lacking in the other or which complements the other. If the two original parents do not provide all the desired characteristics, other sources can be included in the breeding population. In the pedigree method, superior plants are selfed and selected in successive generations. In the succeeding generations the heterozygous condition gives way to homogeneous lines as a result of self-pollination and selection. Typically in the pedigree method of breeding five or more generations of selfing and selection is practiced: F.sub.1 .fwdarw.F.sub.2 ; F.sub.2 .fwdarw.F.sub.3 ; F.sub.3 .fwdarw.F.sub.4 ; F.sub.4 .fwdarw.F.sub.5, etc.\nRecurrent selection breeding, backcrossing for example, can be used to improve an inbred line and a hybrid which is made using those inbreds. Backcrossing can be used to transfer a specific desirable trait from one inbred or source to an inbred that lacks that trait. This can be accomplished, for example, by first crossing a superior inbred (recurrent parent) to a donor inbred (non-recurrent parent), that carries the appropriate gene(s) for the trait in question. The progeny of this cross is then mated back to the superior recurrent parent followed by selection in the resultant progeny for the desired trait to be transferred from the non-recurrent parent. After five or more backcross generations with selection for the desired trait, the progeny will be homozygous for loci controlling the characteristic being transferred, but will be like the superior parent for essentially all other genes. The last backcross generation is then selfed to give pure breeding progeny for the gene(s) being transferred. A hybrid developed from inbreds containing the transferred gene(s) is essentially the same as a hybrid developed from the same inbreds without the transferred gene(s).\nElite inbred lines, that is, pure breeding, homozygous inbred lines, can also be used as starting materials for breeding or source populations from which to develop other inbred lines. These inbred lines derived from elite inbred lines can be developed using the pedigree breeding and recurrent selection breeding methods described earlier. As an example, when backcross breeding is used to create these derived lines in a maize plant breeding program, elite inbreds can be used as a parental line or starting material or source population and can serve as either the donor or recurrent parent.\nDevelopment of Maize Hybrids\nA single cross maize hybrid results from the cross of two inbred lines, each of which has a genotype that complements the genotype of the other. The hybrid progeny of the first generation is designated F.sub.1. In the development of commercial hybrids in a maize plant breeding program, only the F.sub.1 hybrid plants are sought. Preferred F.sub.1 hybrids are more vigorous than their inbred parents. This hybrid vigor, or heterosis, can be manifested in many polygenic traits, including increased vegetative growth and increased yield.\nThe development of a maize hybrid in a maize plant breeding program involves three steps: (1) the selection of plants from various germplasm pools for initial breeding crosses; (2) the selfing of the selected plants from the breeding crosses for several generations to produce a series of inbred lines, which, although different from each other, breed true and are highly uniform; and (3) crossing the selected inbred lines with different inbred lines to produce the hybrid progeny (F.sub.1). During the inbreeding process in maize, the vigor of the lines decreases. Vigor is restored when two different inbred lines are crossed to produce the hybrid progeny (F.sub.1). An important consequence of the homozygosity and homogeneity of the inbred lines is that the hybrid between a defined pair of inbreds will always be the same. Once the inbreds that give a superior hybrid have been identified, the hybrid seed can be reproduced indefinitely as long as the homogeneity of the inbred parents is maintained.\nA single cross hybrid is produced when two inbred lines are crossed to produce the F.sub.1 progeny. A double cross hybrid is produced from four inbred lines crossed in pairs (A.times.B and C.times.D) and then the two F.sub.1 hybrids are crossed again (A.times.B).times.(C.times.D). Much of the hybrid vigor exhibited by F.sub.1 hybrids is lost in the next generation (F.sub.2). Consequently, seed from hybrids is not used for planting stock.\nHybrid seed production requires elimination or inactivation of pollen produced by the female parent. Incomplete removal or inactivation of the pollen provides the potential for self pollination. This inadvertently self pollinated seed may be unintentionally harvested and packaged with hybrid seed.\nOnce the seed is planted, it is possible to identify and select these self pollinated plants. These self pollinated plants will be genetically equivalent to the female inbred line used to produce the hybrid.\nTypically these self pollinated plants can be identified and selected due to their decreased vigor. Female selfs are identified by their less vigorous appearance for vegetative and/or reproductive characteristics, including shorter plant height, small ear size, ear and kernel shape, cob color, or other characteristics.\nIdentification of these self-pollinated lines can also be accomplished through molecular marker analyses. See, \"The Identification of Female Selfs in Hybrid Maize: A Comparison Using Electrophoresis and Morphology\", Smith, J. S. C. and Wych, R. D., Seed Science and Technology 14, pp. 1-8 (1995), the disclosure of which is expressly incorporated herein by reference. Through these technologies, the homozygosity of the self pollinated line can be verified by analyzing allelic composition at various loci along the genome. Those methods allow for rapid identification of the invention disclosed herein. See also, \"Identification of Atypical Plants in Hybrid Maize Seed by Postcontrol and Electrophoresis\" Sarca, V. et al., Probleme de Genetica Teoritica si Aplicata Vol. 20 (1) p. 29-42.\nAs is readily apparent to one skilled in the art, the foregoing are only some of the various ways by which the inbred can be obtained by those looking to use the germplasm. Other means are available, and the above examples are illustrative only.\nMaize is an important and valuable field crop. Thus, a continuing goal of plant breeders is to develop high-yielding maize hybrids that are agronomically sound based on stable inbred lines. The reasons for this goal are obvious: to maximize the amount of grain produced with the inputs used and minimize susceptibility of the crop to pests and environmental stresses. To accomplish this goal, the maize breeder must select and develop superior inbred parental lines for producing hybrids. This requires identification and selection of genetically unique individuals that occur in a segregating population. The segregating population is the result of a combination of crossover events plus the independent assortment of specific combinations of alleles at many gene loci that results in specific genotypes. The probability of selecting any one individual with a specific genotype from a breeding cross is infinitesimal due to the large number of segregating genes and the unlimited recombinations of these genes, some of which may be closely linked. However, the genetic variation among individual progeny of a breeding cross allows for the identification of rare and valuable new genotypes. These new genotypes are neither predictable nor incremental in value, but rather the result of manifested genetic variation combined with selection methods, environments and the actions of the breeder.\nThus, even if the entire genotypes of the parents of the breeding cross were characterized and a desired genotype known, only a few if any individuals having the desired genotype may be found in a large segregating F.sub.2 population. Typically, however, neither the genotypes of the breeding cross parents nor the desired genotype to be selected is known in any detail. In addition to the preceding problem, it is not known how the genotype would react with the environment. This genotype by environment interaction is an important, yet unpredictable, factor in plant breeding. A breeder of ordinary skill in the art cannot predict the genotype, how that genotype will interact with various climatic conditions or the resulting phenotypes of the developing lines, except perhaps in a very broad and general fashion. A breeder of ordinary skill in the art would also be unable to recreate the same line twice from the very same original parents as the breeder is unable to direct how the genomes combine or how they will interact with the environmental conditions. This unpredictability results in the expenditure of large amounts of research resources in the development of a superior new maize inbred line."}
{"text": "The present invention relates to replacing an existing main with a replacement main.\nGas, water and sewage mains may need to be replaced when they develop leaks or for other reasons. Leaks may develop at joints as a result of leak paths appearing within the joints in the case of cast iron gas mains or as a result of cracks in the wall due to corrosion or mechanical damage in cast iron and other materials e.g. ductile iron and steel, cement, brickwork or plastics.\nIn the case where the main to be replaced is a branch main joined to a fluid-carrying spine main, present techniques of mains replacement necessitate an excavation being made at the junction between the branch main and the spine main to gain access thereto. This is frequently in the middle of a busy thoroughfare and the excavation may lead to disruption of both pedestrians and traffic."}
{"text": "Various battery operated hand portable lamps have been developed which provide a source of light where conventional electric power is not available or practical. Such devices are commonly known as flashlights, electric lanterns, or portable utility lights. In addition to requiring lighting in remote places it is also often desirable to be able to place the light source on a flat surface and direct the light beam to most effectively serve the user. Such applications include playing board games at night while camping or working on one's automobile.\nNumerous hand portable lamp devices are presently being marketed including many types of flashlights using one or more size A, size C or size D batteries. Some hand portable lamps use one or more 6 volt dry cells. One such device is the EVERREADY.TM., Model No. 5209 electric lantern known as the \"Indoor/Outdoor Commander\" and available from Union Carbide Corp. of Danbury, Conn. This device uses two 6 volt dry cell batteries and a tubular fluorescent lamp. This lantern, like most flashlights, is essentially comprised of a single body portion and does not allow for independent movement between the lantern body and its lamp. Consequently, if no convenient surface is available to place the lantern on, the user must hold it in order to maximize the lighting in a desired area. The Ray-o-Vac Company markets a portable lantern similar to the EVEREADY Commander in that they both house the illuminescent member and the batteries within a common housing. This type of arrangement, while satisfactory in applications where light positioning is not critical, suffers from the same drawbacks as the EVEREADY Commander and consequently provides undesirable results when the light must be critically positioned.\nOne lantern device which does allow for the positioning of the illuminescent member independent from the lantern's base is shown in U.S. Pat. No. 3,519,812, issued to Brindley and entitled \"Portable Electric Lantern with Swivel Type Lamp Assembly.\" The Brindley patent discloses a lantern case and a swivel type lamp assembly adapted to be elevated above the case for illuminating the surface on which the lantern is placed. Although the Brindley patent does disclose a portable lantern with a swivel lamp, it is not designed to be used in applications where the light emitted from the lantern must be directed to areas other than the surface on which the lantern is placed. The Brindley lantern uses one lamp post to support the overhead lamp assembly. In order to adequately support the assembly, this single post must be quite wide. Thus, another drawback of the Brindley device is the undesirable shadow that is cast on account of its wide lamp post.\nAccordingly, it is a primary object of the present invention to provide a hand portable light with a base housing and extendable lamp housing such that the light emitted from the lamp housing is capable of being positioned in a multiplicity of directions without moving the base housing.\nA further object of this invention is to provide a hand portable light as described above wherein the shadows cast from the extensions which support the lamp housing are minimized, so as to facilitate night-time activities requiring fairly uniform lighting from a portable battery-powered light source.\nStill another object of the invention is to provide a hand portable light having telescoping rod antennas which can support or hold a lamp housing at anyone of a variety of attitudes, as desired by the user."}
{"text": "For some computational tasks, quantum computing devices may offer a computational speed up compared to classical devices. For example, quantum computers may achieve a speed up for tasks such as database search, evaluating NAND trees, integer factorization or the simulation of quantum many-body systems.\nAs another example, adiabatic quantum annealers may achieve a computational speed up compared to classical annealers for some optimization tasks. To perform an optimization task, quantum hardware may be constructed and programmed to encode the solution to a corresponding optimization task into an energy spectrum of a many-body quantum Hamiltonian characterizing the quantum hardware. For example, the solution is encoded in the ground state of the Hamiltonian."}
{"text": "Semiconductor devices are subjected to a series of test procedures in order to confirm functionality and yield, and to assure quality and reliability. This testing procedure conventionally includes \"probe testing\", in which individual dice, while still on a wafer, are initially tested to determine functionality and speed. Probe cards are used to electrically test die at that level. The electrical connection interfaces with only a single die at a time in wafer, before the dice are singulated from wafer.\nIf the wafer has a yield of functional dice which indicates that quality of the functional dice is likely to be good, each individual die is traditionally assembled in a package to form a semiconductor device. Conventionally, the packaging includes a lead frame and a plastic or ceramic housing.\nThe packaged devices are then subjected to another series of tests, which include burn-in and discrete testing. Discrete testing permits the devices to be tested for speed and for errors which may occur after assembly and after burn-in. Burn-in accelerates failure mechanisms by electrically exercising the devices (devices under test or DUT) at elevated temperatures, elevated dynamic biasing schemes. This induces infant mortality failure mechanisms and elicit potential failures which would not otherwise be apparent at nominal test conditions.\nVariations on these procedures permit devices assembled onto circuit arrangements, such as memory boards, to be burned-in, along with the memory board in order to assure reliability of the circuit board and the circuit board assembly and manufacturing process, as populated with devices. This closed assembly testing assumes that the devices are discretely packaged in order that it can then be performed more readily.\nSemiconductor packaging has been referred to in terms of \"levels\" of packaging. The chip capsule generally constitutes a first level of packaging. A second level would then be a \"card\" or a printed circuit board. A third level may include second level packaging combined with a motherboard. A fourth level may follow the third level. In each case, the packaging to any level involves cost.\nIt is proposed that devices be packaged without conventional lead frames. This creates two problems for being conventional test methods. Firstly, discrete testing is more difficult because the conventional lead frame package is not used. Furthermore, multiple device may be assembled into a single package, thereby reducing the performance of the package to that of the die with the lowest performance. This is because the ability to presort the individual dice is limited that obtained through probe testing. Secondly, the packaging may have other limitations package assembly defect mechanisms which are aggravated by burn-in stress conditions so that the packaging becomes a limitation for burn-in testing.\nAccording to the invention represented by U.S. Pat. No. 4,899,107, to Alan Wood and Tim Corbett, a reusable burn-in/test fixture for discrete die is provided. The fixture consists of two halves, one of which is a die cavity plate for receiving semiconductor dice as the devices under test (DUT); and the other half establishes electrical contact with the dice and with a burn-in oven.\nThe first half of the test fixture contains cavities in which die are inserted circuit side up. The die will rest on a floating platform. A support mechanism under the die platform will provide a constant uniform pressure or force to maintain adequate electrical contact to the die contacts on the DUT to probe tips on the second half. The support mechanism will compensate for variations of overall die thickness.\nThe second half has a rigid high temperature rated substrate, on which are mounted probes for each corresponding die pad. Each probe is connected to an electrical trace on the substrate (similar to a P.C. board) so that each die pad of each die is electrically isolated from one another for high speed functional testing purposes. The probe tips are planar so that contact to each die pad occurs simultaneously. The probe tips are arranged in an array to accommodate eight or more dice. The traces from the probes terminate in edge fingers to accept a conventional card edge connector. The geometry of the probes and edge fingers are optimized to avoid electrical test artifacts.\nThe two halves of the test fixture are joined so that each pad on each die aligns with a corresponding electrical contact. The test fixture is configured to house groups of 8 or 16 die for maximum through put efficiency of the functional testers. The test fixture need not be opened until the burn-in and electrical test are completed. After burn-in stress and electrical test, the die are removed from the test fixture and depositioned accordingly. The fully burned-in and tested die are available for any type of subsequent assembly applications.\nThat technique allows all elements of the burn-in/test fixture to 100% reusable, while permitting testing of individual dice in a manner similar to that accomplished with discrete packaged semiconductor devices.\nAn ability to extend accelerated burn-in and functional/parametric/speed testing of dice to include accelerated burn-in and functional, parametric and speed testing while the dice are still on the wafer would have several advantages. Since each step in the assembly and package process represents commitment of resources, early determination of defective parts or ability to predict a failure at a conventional burn-in stage is advantageous. It would be further advantageous to be able to predict a failure at a burn-in stage prior to assembly. Clearly, if a part can be made to fail prior to assembly, assembly resources can be directed to a higher percentage of good parts.\nThere exists a significant market for uncut fabricated wafers. These wafers are referred to as \"probe wafers\" because they are delivered after probe testing, which follows fabrication. The purchase of probe wafers is primarily by \"ASIC assembly houses\" which custom package parts, including parts traditionally considered to be \"commodity\" chips. The purchase of uncut wafers is usually based on the recent yield rate of the semiconductor manufacturer, but recent yields are not a strong indicator of the yield of any given wafer lot. Furthermore, the assembly process techniques used by the assembly house have a significant effect on yield.\nCharacterization, such as speed grading is even more variable than yield. While a packaged DRAM is purchased by the consumer based on the part's speed grade, speed grading of probe wafers is almost a matter of conjecture. That means that it is happenchance as to whether the assembly house purchases a wafer of mostly \"-10\" parts (100 ns) or mostly \"-6\" parts (60 ns).\nRecent developments in fabrication technology have resulted in such speed characterizations being more uniform on any given wafer. This has made it possible to provide wafers in which a majority of good dice have speed grades which do not greatly exceed an average for the wafer. Such uniformity, along with an ability to achieve fuse repairs and patches, have made wafer scale integration of arrays and cluster packaging practical.\nOther developments include an ability to track individual dice on wafers, starting at probe. Traditionally, probe identifies bad dice, (example, an ink spot.) The assembly process is continued only for dice which do not have the ink spots. By computer tracking, the ink spot becomes superfluous, as a map of good and bad dice are stored and transferred to subsequent assembly steps.\nAlthough the dice are singulated, there are cases in which the discrete parts are reassembled into an array after assembly. An example is in computer memory, in which one or more banks of memory are composed of arrays of memory chips. It would be advantageous to be able to select good dice on a wafer and assemble the dice into an array without singulating the dice. This would allow a much denser array of good clustered dice on a single piece of silicon.\nIt is an object of the invention to increase handling efficiency, while at the same time reducing the required size of the test fixture."}
{"text": "Teeth are lost for a variety reasons, such as tooth decay, root canal failure, periodontitis (gum disease), trauma to the mouth, excessive wear and tear, and congenital defects. Patients who have lost teeth are often reluctant to smile or talk. Additionally, biting irregularities caused by tooth loss can have a negative effect on eating habits. This can lead to secondary health problems, such as malnutrition.\nRegardless of the nature of problems related to tooth loss, dental implants may provide a simple remedy with proven results. Dental implants are prosthetic (i.e., artificial) tooth replacements that are used to counter tooth loss. Dental implants are stronger and more durable than their restorative counterparts (i.e., dentures), and offer a permanent solution to tooth loss. Additionally, implants may be used in conjunction with other restorative procedures for maximum effectiveness. For example, a single implant can serve to support a crown replacing a single missing tooth. Implants can also be used to support a dental bridge for the replacement of multiple missing teeth, and can be used with dentures to increase stability and reduce gum tissue irritation.\nDental implant stents and other devices presently used by surgeons to place dental implants in a patient's mouth do not have features that allow the surgeons to place a dental implant parallel to other implants and/or existing teeth and at precise, predetermined depths within the arch of the mandible or maxilla. A similar drawback occurs with dental implants used as supporting members for the development of a dental prosthesis. When dental implants are not placed parallel to other dental implants, final restoration prostheses are much more difficult to manufacture and the fit and longevity of a final restoration is reduced.\nPresent surgical dental implant stents also lack features that allow surgeons to place a dental implant at a predetermined depth decided upon by the surgeon or allow a dental lab technician to manufacture a surgical dental implant stent that can be used to fabricate the final dental implant restoration that will be reliably retained in the patient's mouth. When surgeons have no means to accurately form an aperture to place a dental implant at a predetermined depth, the surgeons when making this aperture may pass through the bone of the mandible or maxilla, necessitating the location of a new site as well as restoring the damaged area. This often increases the amount of time required for the patient to heal and recover from the procedure.\nCurrent surgical dental implant stents also lack features that would allow a surgeon to immediately place a final dental implant restoration. In fact, current dental implant stents normally take weeks to complete due to the trial-and-error fabrication process inherent to present systems. There is a need for a way to more accurately position oral implants in a patient's mouth and to reduce the amount of time required to produce and install them."}
{"text": "The invention relates to an apparatus for conditioning of bulk material, such as shavings or chips, especially of the vegetable type, such as wood-chips, for the production of fibers or pulp for making paper or cardboard or making fiberboards.\nIt is known to produce the pulp necessary for paper-making by a process in which wood chips or other vegetable substances are used. These wood chips, which can be approximately 20 to 40 mm long and up to approximately 20 mm wide, are defibered, so as then to be processed further in the appropriate way. To achieve as high a fiber quality as possible, the wood chips can be conditioned in a silo by means of low-pressure or high pressure steam.\nThe conditioned bulk material is discharged by means of a clearing unit through one or more orifices provided in the silo bottom, via screw conveyors, and is delivered to the defibration mill. During interruptions in operation, for example, during operating faults and the like, blockages in the silo occur as a result of the swelling of the wood-chips. Such agglomerations can be eliminated only by knocking, poking, etc. Moreover, the bulk material tends to form bridges, with a result that the discharge of the material from the silo becomes irregular and can even be blocked completely. Also, the temperature of the wood-chips is uneven and is unsatisfactory in view of the subsequent process steps. Considerable steam losses occur as a result of the formation of bridges and channels in the bulk material. Another disadvantage is that, because of the irregular steam distribution or varying retention time, the wood-chips change irregularly and their properties are thereby impaired, so that during subsequent use problems of quality can arise and additional costs are incurred."}
{"text": "The present invention relates to the testing of hydraulic fluid systems.\nThe satisfactory operation of hydraulic fluid control systems depends on the integrity of their sealing components and the failure of even one such component can seriously impair the operating characteristics of the hydraulic system. In particular, the failure of one or more sealing components can result in a significant reduction in system pressure, preventing proper operation of the devices, such as valves, which are controlled by the system.\nCurrent practice in the industry is to perform scheduled preventive maintenance of critical parts. Failures which do not prevent operation of the system are usually ignored until an operating failure occurs, at which time suspected components of the hydraulic fluid system are replaced one at a time and after each replacement, the installation is placed back into operation to determine whether the replaced component is the one responsible for the system failure.\nIn any event, each failure of the hydraulic control system requires shutdown of the installation, with attendant lose of productivity.\nMany installations of the type hereunder consideration must undergo periodic shutdown for various reasons unrelated to the operating state of the hydraulic control system. If, during these shutdown periods, components of the hydraulic control system which have begun to fail but continue to be at least marginally operative were located, then these components could be replaced to thereby avoid the need for subsequent, unscheduled shutdowns."}
{"text": "Heretofore various types of coin and currency changers have been known and utilized. However, most known changers have certain drawbacks with regards to the jamming of coin chutes, the necessity of stacking the coins in a magazine, non-positive payouts, jackpotting, and the inflexibility of the system to adapt to payouts of varying combinations of change for any given currency or coin.\nConsequently, it is an object of the instant invention to present a changer utilizing an angled edge stacking coin chute which eliminates the inherent problem of jamming due to coin deformation or the insertion of an improper denomination of coins into a given chute.\nYet another object of the invention is to present a changer which utilizes a bulk hopper wherein coins may be loaded without stacking and from which coins may be uniquely transported for depositing into the angled edge stacking chute.\nStill another object of the invention is to present a changer which utilizes a positive payout technique wherein accounting is made of the number of coins actually dispensed rather than the number of times which a coin dispensing mechanism has functioned.\nStill another object of the invention is to present a changer utilizing an overpayment detector circuit which immediately terminates the dispensing cycle and inhibits further operation of the changer when a first coin, constituting an overpayment, is dispensed.\nA further object is to provide an instantaneously data controlled change system for providing variable change combinations or amounts.\nYet a further object of the invention is to present a changer which is field programmable by a serviceman to preselect the particular combination of coins to be dispensed for any particular change-making operation so as to optimize the utility of the changer.\nThese objects and other objects which will become apparent as the detailed descripton proceeds are achieved by a positive payout coin and currency changer operative to dispense change from money of any of numerous denominations, comprising: first circuit means for determining the validity of the money for which change is desired and producing a signal indicative of the validity and denomination thereof; coin chutes, one for each denomination of coins to be dispensed, operative to receive coins from associated hoppers and maintain said coins therein, and second circuit means interconnecting said first circuit means and said chutes for receiving and responding to said signal to effectuate the release of respective predetermined numbers of coins from each chute and operative to sense the escapment of each coin from each chute."}
{"text": "1. Field of the Invention\nThis invention relates to the field of medical products. More specifically, the invention comprises lumbar belt which allows for easier application and adjustment. The lumbar belt may be used in a variety of therapies, including cryotherapy.\n2. Description of the Related Art\nLumbar belts are used as part of a variety of therapies. One example is the use of a “cold therapy” (also known as “cryo therapy”) to minimize swelling following a surgical procedure. Cryo therapy assumes a variety of forms—ranging from simple ice or gel packs to sophisticated chilled water circulating devices.\nOne common application for cryo therapy is in post-surgical swelling reduction for lumbar spinal surgery. It is known to use a compressive “belt” to hold cold packs in position over the person's lumbar spine. The “belt” needs to be properly positioned and easily adjustable for the user. The present invention provides features and methods to facilitate the installation and provide for quick adjustments for such a belt."}
{"text": "1. Field of the Invention\nThis invention relates to isolating the code and/or data of one or more software modules within a computer system.\n2. Description of the Related Art\nThe invention may be implemented as one or more computer programs or as one or more computer program modules embodied in one or more computer readable media. The computer readable media may be based on any existing or subsequently developed technology for embodying computer programs in a manner that enables them to be read by a computer. For example, the computer readable media may comprise one or more CDs (Compact Discs), one or more DVDs (Digital Versatile Discs), some form of flash memory device, a computer hard disk and/or some form of internal computer memory, to name just a few examples. An embodiment of the invention, in which one or more computer program modules is embodied in one or more computer readable media, may be made by writing the computer program modules to any combination of one or more computer readable media. Such an embodiment of the invention may be sold by enabling a customer to obtain a copy of the computer program modules in one or more computer readable media, regardless of the manner in which the customer obtains the copy of the computer program modules. Thus, for example, a computer program implementing the invention may be purchased electronically over the Internet and downloaded directly from a vendor's web server to the purchaser's computer, without any transference of any computer readable media. In such a case, writing the computer program to a hard disk of the web server to make it available over the Internet may be considered a making of the invention on the part of the vendor, and the purchase and download of the computer program by a customer may be considered a sale of the invention by the vendor, as well as a making of the invention by the customer.\nThis invention may be implemented in a wide variety of computer systems, having a wide variety of hardware platforms and configurations and a wide variety of software platforms and configurations. If a computer system includes multiple software entities or software modules, or at least has the potential to contain multiple software modules, then the integrity of the computer system may be improved by implementing this invention to protect the code and/or data of one or more of the software modules from other software modules in the system.\nOver the years, a variety of techniques have been used for executing multiple software modules within a computer system. Early computer systems could execute multiple software programs, but they could only execute one program at a time. Such computers might load one program into memory and execute it to completion or other termination, before proceeding to a subsequent program that would then be loaded into memory and executed. As another example, various multitasking operating systems (OSs) enable multiple programs (or selected portions thereof) to be loaded into memory at one time and executed in an alternating manner, according to a scheduling algorithm. Also, some processors include multithreading capabilities, which enable multiple threads of one or more programs to be executed simultaneously on a single processor. Finally, multiprocessor computer systems have also become commonplace, in which each of multiple processors can execute one or more threads all at the same time. This invention may be advantageously implemented in any of these types of systems, as well as other possible computer systems in which multiple software modules may be executed.\nSuch computer systems generally attempt to isolate the code and data of the different software modules within the computer system from each other, so that, for example, one software module cannot interfere with the execution of another software module by altering its code or data. Such isolation may be provided for code and/or data that is stored on a hard drive (or other secondary data storage means) and/or that is resident in main memory (or other primary data storage means). In this patent, the term “data” is generally used in a broad sense, to include data that is operated on by the instructions (code) of a software module as well as the contents of a stack and any other possible forms of data that are associated with a software module. As one example of the isolation of code and data, many systems implement a virtual addressing mechanism, in which different software modules within the system have different virtual address spaces, with each virtual address space generally being mapped to different portions of the physical address space of the computer system, so that the virtual addresses of a given software module are generally only mapped to physical addresses that contain the code or data of that particular software module. Virtual addressing mechanisms are described in greater detail below. A given software module may attempt to access every memory location in its own virtual address space, accessing every memory location to which it has access, and it will still only be able to access its own code and data (assuming that there is no shared memory). Thus, providing a virtual addressing mechanism provides some isolation between the code and data of multiple software modules in a computer system. Various other protection mechanisms may also be implemented in such computer systems to isolate the code and/or data of multiple software modules from one another.\nAlthough the invention may be implemented in a wide variety of computer systems, having a wide variety of hardware and software platforms and configurations, the following description is generally limited to a single hardware platform for brevity. In particular, this description is generally limited to computer systems that include one or more processors having the “x86” architecture, which is described in the IA-32 Intel Architecture Software Developer's Manual (“the IA-32 Manual”) from Intel Corporation. Also for brevity, the following description is generally limited to computer systems running a Windows OS from Microsoft Corp. or a Linux OS, although there are certainly other OSs that operate on the x86 platform. A Windows OS from Microsoft Corp. may be a Windows XP OS or a Windows 2000 OS, for example, while a Linux OS may be a distribution from Novell, Inc. (SUSE Linux), Mandrakesoft S.A. or Red Hat, Inc. Based on the following description related to the x86 architecture and a Windows or Linux OS, a person of skill in the art will be able to implement the invention in a wide variety of other computer systems.\nThe x86 architecture provides two primary memory protection mechanisms that may be used by an OS (or other system software) to try to isolate the code and data of multiple tasks or processes that execute on the processor, namely a segmentation mechanism and a paging mechanism. The IA-32 Manual may be consulted for a detailed description of these protection mechanisms. The Windows and Linux OSs use the paging mechanism, but they generally don't take advantage of the segmentation mechanism. Instead, these OSs define segments that include the entire addressable range of the processor, so that the segmentation protection mechanism becomes ineffective in providing isolation between the code and data of multiple tasks. Thus, for simplicity, this discussion focuses on the paging mechanism of the x86 processor, which implements a virtual addressing mechanism as described briefly above and in greater detail below.\nVery briefly for now, for Windows and Linux OSs, different user processes are generally given different virtual address spaces. The OS creates a different set of page tables (and a page directory) for each virtual address space, which maps the respective virtual addresses to physical addresses. Thus, the page tables for a given user process map that process's virtual addresses to the physical addresses that contain the code and data for that process. The page tables for the user processes also contain mappings for code and data of the OS, but the user processes cannot use these mappings because the user processes are executed at a Current Privilege Level (CPL) of 3 and these mappings are set to require a supervisor privilege level (a CPL of 0, 1 or 2). Otherwise, the page tables for a given user process generally only contain mappings to physical memory pages that contain that process's code and data. Therefore, a user process can generally only access its own code and data. Executing the user processes at a CPL of 3, also prevents the processes from modifying their own page tables. Otherwise, a process could add entries to its page tables that map to any physical address in the system, so that the process could give itself access to the code and data of other software modules, including other user processes and the OS.\nWindows and Linux OSs generally provide adequate protection for the software modules in a computer system, so long as all of the software modules are well designed and well behaved, meaning that they are not attempting to circumvent the protection mechanism. Thus, many processes may be running concurrently in such a computer system, with the OS giving each process a share of the system resources, including processor time, memory space and hard disk space, without any of the processes interfering with the code or data of the other processes.\nFIG. 1 illustrates a simple computer system 2A having multiple software modules. The computer system 2A includes system hardware 100A, an OS 20A, a first application 40A and a second application 40B. The system hardware 100A may be conventional hardware based on, for example, the x86 platform, and the OS 20A may be, for example, a Windows OS or a Linux OS. The applications 40A and 40B may be any applications designed to run on the system hardware 100A and the OS 20A. The OS 20A also includes a set of drivers 29A, which may be conventional drivers for the OS 20A, possibly including one or more drivers from a company that is different from the OS vendor (a third party vendor).\nThe OS 20A, in conjunction with the system hardware 100A, attempts to isolate the code and data of the applications 40A and 40B from one another. For example, the OS 20A and the system hardware 100A may implement a virtual addressing mechanism, as described above. As illustrated in FIG. 1, implementing such a protection mechanism may be characterized as establishing an isolation barrier 80B between the applications 40A and 40B, preventing (or at least hindering) one application from accessing the code and data of the other application. There may also be some code and/or data that is shared explicitly or transparently between the applications 40A and 40B. Techniques are known for allowing such sharing of code and data, while maintaining isolation between the applications 40A and 40B. For example, the OS 20A may mark physical memory pages that contain shared code or data as read only, such as when using a copy-on-write (COW) technique. The isolation barrier 80B may be referred to as an “OS isolation barrier” because it is implemented by the OS 20A, in conjunction with the system hardware 100A. The OS 20A, again in conjunction with the system hardware 100A, also establishes an OS isolation barrier 80A between the OS 20A and all applications in the system, including the applications 40A and 40B, so that the applications are prevented (or hindered) from directly accessing the code and data of the OS 20A. In the case of a Windows or Linux OS running on an x86 platform, the OS isolation barrier 80A is established by executing the applications in the system at a CPL of 3 and requiring a supervisor privilege level to access memory pages containing the code and data of the OS 20A.\nAlthough the Windows and Linux OSs provide adequate isolation between software modules for computer systems that contain only well designed and well behaved software modules, malicious software modules have been known to wreak havoc in such computer systems by circumventing the protection mechanisms and engaging in all sorts of mischief. In particular, such malicious software modules have been known to breach the OS isolation barriers 80B and 80A, and corrupt the code and/or data of other applications in the system, and/or of the OS itself. Numerous security vulnerabilities have been discovered in the Windows OSs and, to a lesser extent, in the Linux distributions, and many of these vulnerabilities have been exploited by hackers using different types of malicious software, such as viruses, worms, etc. Microsoft Corp. acknowledges in its '006 patent (referenced below) that “open operating systems that allow users to easily install hardware and software . . . are inherently untrustworthy” (the '006 patent, paragraph [0004]). Poorly designed or implemented software may inadvertently bypass these protection mechanisms too and may also wreak havoc in a computer system. Although the description in this patent generally relates to malicious software, it also applies to software that inadvertently has the same or similar effects as malicious software.\nThus, hackers that exploit the vulnerabilities described above do so for a variety of reasons and with a variety of goals, some being relatively benign and others being quite destructive or disruptive. As one example, a malicious software module may be written and deployed that searches for sensitive data on a computer's hard drive or in its memory and transmits any such sensitive data back to the hacker that launched the malicious code. If the malicious software manages to execute at supervisor level (with a CPL of 0, 1 or 2), which is a common occurrence, then the malicious software may access a variety of system resources that are supposed to be safeguarded. For example, the malicious software can create its own page table entries, and obtain access to any address within the physical address space of the processor. The malicious software can then scan the entire physical memory for sensitive data, including memory that contains code and data of other software modules. A wide variety of other possibilities also exist.\nSecurity threats such as these have been gaining greater notoriety, and it is widely accepted that something should be done to improve the security of the ubiquitous personal computer, and, in particular, there is a recognized need to improve the security for the vast number of computers based on the x86 architecture. Many people believe that software changes alone will not provide adequate protection. Accordingly, many different companies are working toward solutions that involve substantial changes to both the system hardware and the system software of a computer system. Many such security measures will also require substantial changes to application level software as well. Some of the security measures that have been or are being pursued by various companies are described below.\nU.S. Pat. No. 6,507,904 (Ellison et al., “Executing Isolated Mode Instructions in a Secure System Running in Privilege Rings”, “the '904 patent”), which was assigned to Intel Corporation, describes a computer system in which a processor operates in either a normal execution mode or an isolated execution mode. Among other features, the computer system includes an isolated region in system memory that is protected by both the processor and a chipset of the computer system. Access to the isolated region is permitted only when using special bus cycles, referred to as isolated read and write cycles. The isolated read and write cycles may only be issued by the processor when it is executing in the isolated execution mode. Execution in the isolated execution mode is restricted, so that access to the isolated memory region is also restricted, accordingly.\nU.S. Pat. No. 6,820,177 (Poisner, “Protected Configuration Space in a Protected Environment”, “the '177 patent”), which was also assigned to Intel Corporation, describes a computer system in which the physical address space encompasses a protected configuration space and a non-protected configuration space. Attempts to access addresses within the protected configuration space are redirected to protected configuration hardware that is external to the system memory. This protected configuration hardware holds control values, provides control information, and performs operations pertaining to a protected operating environment. Attempts to access the protected configuration space are permitted only if they are made by a processor, and only if they are made using a “protected command.” Restricting access to the protected configuration hardware purportedly enables the system to maintain a protected operating environment. The protected operating environment may include blocks of protected memory, which are apparently protected from any attempted access by non-processor devices.\nU.S. Pat. No. 6,986,006 (Willman et al., “Page Granular Curtained Memory via Mapping Control”, “the '006 patent”), which was assigned to Microsoft Corporation, describes a method by which access to trusted memory is restricted using a paging mechanism, by not including mapping entries in page tables that map to physical memory pages that contain the trusted memory. The memory pages that contain the page tables are then restricted to read-only access when the processor is operating in a non-trusted mode to prevent non-trusted software from adding a new mapping entry or modifying an existing mapping entry to map to trusted memory. If non-trusted software attempts to write to a memory page containing a page table, a context switch is initiated into a page table entry edit module, which is trusted software. The page table entry edit module then ensures that the attempted write does not establish a mapping into trusted memory. The '006 patent does not specifically indicate how non-trusted software is prevented from changing the memory pages containing page tables to read/write access or from corrupting or supplanting the page table entry edit module to enable the non-trusted software to establish a mapping to trusted memory. The '006 patent does indicate, however, that the memory controller or other hardware may be able to restrict access to certain pages of physical memory under the control of the page table entry edit module. A new hardware platform or substantial changes to an existing hardware platform would presumably be necessary to implement such a restriction.\nU.S. Pat. No. 7,058,768 (Willman et al., “Memory Isolation Through Address Translation Data Edit Control”, “the '768 patent”), which is a continuation-in-part of the '006 patent and which was also assigned to Microsoft Corporation, is similar to the '006 patent, although the '768 patent uses different terminology and it goes into greater detail in some areas. The '768 patent describes a computer system that includes a trusted environment and an untrusted environment. A trusted component in the trusted environment purportedly ensures that, when the system is executing in the untrusted environment, no active address translation map includes an address mapping that maps to isolated memory, so that the isolated memory is not accessible to untrusted components. Again, the '768 patent does not specify how the trusted environment is established or maintained, but implementing the computer system described in the '768 patent presumably requires substantial hardware changes.\nU.S. Patent Application Publication No. 2004/0205203 (Peinado et al., “Enforcing Isolation Among Plural Operating Systems”, “the '203 application”), which was also assigned to Microsoft Corporation, describes a method for restricting the physical addresses that are accessible to Direct Memory Access (DMA) devices in a computer system in which multiple OSs run. In this method, a security kernel maintains a DMA exclusion vector that specifies which physical addresses are accessible to different DMA devices. A hardware device, referred to as a regulator, enforces the physical address restrictions specified in the DMA exclusion vector. The '203 application also briefly mentions the possibility of using a shadow page table technique to prevent one OS from accessing the private data of another OS, and the '203 application mentions the more general possibility of employing an “adjunct memory access control scheme.”\nU.S. Pat. No. 6,651,171 (England et al., “Secure Execution of Program Code”, “the '171 patent”), which was also assigned to Microsoft Corporation, describes a system in which hardware enforces a restricted memory arrangement. Multiple curtained memory rings are arranged in a hierarchical manner, similar to the protection rings of the x86 architecture. Different code and data are associated with each of the memory rings. Software in more-privileged rings can access code and data in less-privileged rings, but software in less-privileged rings cannot access code or data in more-privileged rings. Also, there may be multiple subrings within a given ring. Software within a given ring can also access code and data within its own ring, except that it cannot access code or data within a different subring.\nU.S. Patent Application Publication No. 2003/0093686 (Barnes et al., “Memory Management System and Method Providing Linear Address Based Memory Access Security”, “the '686 application”), which was assigned to Advanced Micro Devices, Inc., describes a Memory Management Unit (MMU) that includes a Security Check Unit (SCU) that receives a linear address generated during the execution of a current instruction. The linear address has a corresponding physical address residing within a selected memory page. The SCU uses the linear address to access one or more security attribute data structures located in the memory to obtain a security attribute of the selected memory page. The SCU compares a numerical value conveyed by a security attribute of the current instruction to a numerical value conveyed by the security attribute of the selected memory page, and produces an output signal dependent upon a result of the comparison. The MMU accesses the selected memory page dependent upon the output signal. The security attribute of the selected memory page and the security attribute of the current instruction may each include a security identification (SCID) value indicating a security context level of the selected memory page or the current instruction, respectively.\nU.S. Pat. No. 6,823,433 (Barnes et al., “Memory Management System and Method for Providing Physical Address Based Memory Access Security”, “the '433 patent”), which was also assigned to Advanced Micro Devices, Inc., discloses a MMU that is similar to the MMU disclosed in the '686 application. A security unit in the MMU of the '433 patent uses a physical address, instead of a linear address, to access security attribute data structures to obtain a security attribute of a selected memory page. Otherwise, the disclosure of the '433 patent is substantially similar to the disclosure of the '686 application.\nU.S. Pat. No. 7,073,059 (Worley, J R. et al., “Secure Machine Platform that Interfaces to Operating Systems and Customized Control Programs”, “the '059 patent”), which was assigned to Hewlett Packard Company, describes a “Secure Platform” (SP), which includes a software layer that executes at a privileged level on a modern processor architecture, such as the IA-64 processor architecture from Intel Corporation. The SP interfaces with one or more OSs and customized control programs, and allows them to access non-privileged machine instructions and registers. However, the OSs and customized control programs purportedly have no direct access to privileged instructions and registers and firmware interfaces. Instead, the SP allows the OSs and customized control programs to invoke software routines that provide control of the hardware, without exposing the privileged machine instructions and registers. The SP also organizes the resources of the system into one or more disjoint, mutually isolated partitions, called “domains.” A single OS and all user-level processes managed by that OS together comprise a domain, and each domain is allocated a separate portion of virtual memory, along with other system resources. The SP employs the region registers and region identifiers, along with protection keys, of the IA-64 processor architecture to partition memory between the different domains.\nWith respect to the ubiquitous x86 platform, each of the possible security measures described above would require substantial hardware changes or an entirely new hardware platform. They would also require substantial changes to existing software platforms, including system software and possibly application software. Applications in some of these implementations might also have limited access to input/output devices because of a limited supply of trusted device drivers. A widespread transition from the x86 platform to a new hardware platform will likely be a slow, gradual process. The amount of money that is invested in computer hardware and software based on the x86 architecture throughout the world is enormous. Many individuals, businesses, schools, governments and other organizations will be reluctant to scrap their current x86 systems, along with all the software that currently runs on x86 systems, and replace them with new technology. Even if a new, more secure and widely accepted hardware platform were available today, it would still take a long time for the new hardware to become anywhere near as widespread as the x86 platform is today. In the meantime, a large number and proportion of computers would remain vulnerable to the security threats described above.\nAlso, even those individuals and organizations that quickly adopt the new hardware and software technology may still be susceptible to adverse effects resulting from lingering vulnerabilities of x86 computers. For example, a malicious software module may attack an x86-based server computer running a web application. If the malicious software takes down the server computer, any client computer attempting to access the web application will be adversely affected no matter how secure the client computer is. As another example, a malicious software module may infect a large number of x86 computers connected to the Internet. The many instances of the malicious software module on all the x86 computers may then mount a coordinated denial of service attack on a particular website on the Internet. In this case, the server computer(s) that are hosting the website and any bonafide client computers trying to access the website may be adversely affected by the attack, again, regardless of the technology implemented in either the server computer(s) or the bonafide client computers.\nAs a final example, even if an individual or organization has spent the money to replace all of its own computers and software with the new technology, sensitive information of the individual or organization may still reside on some external computer that remains vulnerable to the security threats described above. Suppose, for example, that an individual forks out the money to buy a new computer system based on a new hardware and software platform, so that all the sensitive information on the individual's computer system is now relatively secure. Suppose further, however, that the individual's bank (or any other organization that has sensitive information of the individual, such as some small company running an obscure Internet website from which the individual has made a purchase) has thus far continued to use its x86 computers instead of investing in the new technology. Again, if malicious software is able to exploit a vulnerability in the bank's computer system, a hacker may be able to steal sensitive information of the individual no matter what technology the individual is using.\nWhat is needed, therefore, is a security measure that can be implemented more quickly and easily, without requiring such a large investment in new computer hardware and software. More particularly, what is needed is a solution that provides better isolation for the code and/or data of software modules within a computer system, but that can be implemented in software, without any hardware changes and with, at most, only minor software changes.\nThere are some proposed security measures that may be implemented primarily in software. In particular, there are some such measures that use virtualization technology to create multiple virtual machines (VMs), where different software modules run in different VMs. It is widely recognized that a well-designed and implemented virtualization layer can generally provide much greater isolation between multiple VMs than a general OS can provide between multiple software modules. For example, U.S. Pat. No. 6,922,774 (Meushaw et al., “Device for and Method of Secure Computing Using Virtual Machines”, “the '774 patent”), which was assigned to the National Security Agency of the U.S. Government, describes a computer system that includes a Virtual Machine Monitor (VMM) that allows a user to create a number of VMs. The system includes a user-definable number of non-sensitive VMs and a user-definable number of sensitive VMs, all of which are isolated from one another by the virtualization technology. Each sensitive VM provides access to a secure area in a computer system that is accessible only through encrypted and/or authenticated connections. An encryption VM is created for and connected to each sensitive VM, where each encryption VM provides encryption capabilities, as well as possibly digital signature and key exchange capabilities. The computer system may also include a server connected to each VM in the system, where the server may be another VM or a stand-alone device. Each VM in the system can send information to the server, and the server can send information to any VM in the system, except as limited by user-definable rules for when a transfer is, or is not, appropriate to be transferred from one VM to another. Thus, the server allows information to be transferred from one VM to another when appropriate while maintaining isolation between VMs.\nAnother proposed security measure that takes advantage of the isolation provided by implementing multiple VMs was described in a technical paper entitled “Terra: A Virtual Machine-Based Platform for Trusted Computing,” which was submitted to and presented at the Symposium on Operating Systems Principles, October 19-22, 2003, by Tal Garfinkel, Ben Pfaff, Jim Chow, Mendel Rosenblum and Dan Boneh (“the Terra paper”). The Terra paper describes a trusted VMM that partitions a single tamper-resistant, general-purpose platform into multiple isolated VMs. Existing applications and OSs can each run in a standard (open-box) VM that provides the semantics of today's open platforms. Applications can also run in their own closed-box VMs that provide the functionality of running on a dedicated closed platform. Among various other aspects of the proposed security measure, VMs on a single physical machine communicate with one another over virtualized standard input/output interfaces such as network interface cards, serial ports, etc.\nFIG. 2 illustrates a general computer system 2B in which multiple VMs are implemented to isolate multiple software modules from one another. The computer system 2B includes system hardware 100B, which may be conventional hardware, such as hardware based on the x86 platform. The system hardware 100B may be substantially the same as the system hardware 100A of FIG. 1, or it may be substantially different. A virtualization software 200A executes on the system hardware 100B and supports a plurality of VMs, such as a first VM 300A and a second VM 300B, in a known manner. The virtualization software 200A may comprise a VMM, for example, such as a VMM as implemented in a virtualization product of VMware, Inc., the assignee of this patent. Such a VMM and other possible units of the virtualization software 200A are described in greater detail below.\nIn supporting the VM 300A, the virtualization software 200A virtualizes a virtual system hardware 310A, which may be based on an existing hardware platform, such as the x86 platform. An OS 20B, along with a set of drivers 29B, run on the virtual system hardware 310A. The OS 20B may be any OS designed to run on the hardware platform virtualized in the virtual hardware 310A. For example, if the virtual hardware 310A is based on the x86 platform, the OS 20B may be a Windows OS or a Linux OS. The set of drivers 29B may be conventional drivers for the OS 20B. A first application 40C and a second application 40D run on the OS 20B. The applications 40C and 40D may be any applications designed to run on the platform of the virtual hardware 310A and the OS 20B. Similar to the OS 20A of FIG. 1, the OS 20B, in conjunction with the virtual system hardware 310A, attempts to isolate the code and data of the applications 40C and 40D from one another, establishing an OS isolation barrier 80B between the applications 40C and 40D. Also similar to the OS 20A of FIG. 1, the OS 20B, again in conjunction with the virtual system hardware 310A, also establishes an OS isolation barrier 80A between the OS 20B and all applications in the VM 300A, including the applications 40C and 40D. Thus, the VM 300A may be substantially the same as the computer system 2A, except that the virtual system hardware 310A is virtual hardware, virtualized by the virtualization software 200A, instead of physical hardware.\nIn supporting the VM 300B, the virtualization software 200A virtualizes a virtual system hardware 310B, which may be based on an existing hardware platform, such as the x86 platform. An OS 20C, along with a set of drivers 29C, run on the virtual system hardware 310B. The OS 20C may be any OS designed to run on the hardware platform virtualized in the virtual hardware 310B. For example, if the virtual hardware 310B is based on the x86 platform, the OS 20C may be a Windows OS or a Linux OS. The set of drivers 29C may be conventional drivers for the OS 20C. A first application 40E and a second application 40F run on the OS 20C. The applications 40E and 40F may be any applications designed to run on the platform of the virtual hardware 310B and the OS 20C. Again, similar to the OS 20A of FIG. 1, the OS 20C, in conjunction with the virtual system hardware 310B, attempts to isolate the code and data of the applications 40E and 40F from one another, establishing an OS isolation barrier 80B between the applications 40E and 40F. Also similar to the OS 20A of FIG. 1, the OS 20C, again in conjunction with the virtual system hardware 310B, establishes an OS isolation barrier 80A between the OS 20C and all applications in the VM 300B, including the applications 40E and 40F. Thus, again, the VM 300B may be substantially the same as the computer system 2A, except that the virtual system hardware 310B is virtual hardware, virtualized by the virtualization software 200A, instead of physical hardware.\nThe virtualization software 200A isolates the VMs in the computer system 2B from one another. For example, the virtualization software 200A allows software within the VM 300A to access portions of physical memory in the system hardware 100B and it allows software within the VM 300B to access other portions of the physical memory. The virtualization software 200A maps attempted memory accesses from the respective VMs 300A and 300B to different portions of the physical memory, ensuring that no memory address generated by software in one VM can access code or data of another VM. In a similar manner, the virtualization software 200A maps attempted hard disk accesses from the respective VMs 300A and 300B to different portions of one or more hard disks in the system hardware 100B, ensuring that one VM cannot access the hard disk space of another VM.\nThe virtualization software 200A also takes other precautions to isolate the VMs in the computer system 2B from one another, and from the virtualization software 200A, itself. For example, U.S. patent application Ser. No. 10/917,713 (Agesen et al., “Restricting Memory Access to Protect Data when Sharing a Common Address Space”, “the '713 application”), which has been assigned to the same assignee as this patent, and which is hereby incorporated herein by reference, describes methods that may be used to enable a VMM to occupy a portion of a linear address space of a VM, while preventing the VM from accessing the memory of the VMM. There are also various other methods that may be used to enable virtualization software to coexist with VMs in a virtual computer system, while protecting or isolating the virtualization software from software within the VMs. The virtualization software 200A may also prevent software within the VMs 300A and 300B from directly accessing certain hardware resources to further isolate the VMs from one another and from the virtualization software 200A. For example, the virtualization software 200A may prevent software within the VMs 300A and 300B from directly accessing a Direct Memory Access (DMA) device to prevent the possibility that the DMA device could be used to access either the hard disk space or the memory of other VMs or of the virtualization software itself. Various other precautions may also be taken, depending on the particular implementation. In addition to the precautions described thus far, the underlying system hardware may provide further support for isolating the VMs from each other and from the virtualization software. For example, Intel Corporation has implemented its Virtualization Technology (Intel VT), and Advanced Micro Devices, Inc. has implemented its AMD Virtualization (AMD-V) or Secure Virtual Machine (SVM), which provide hardware enhancements to facilitate the implementation of virtual computer systems.\nThus, the virtualization software 200A, in conjunction with the system hardware 100B, may be said to establish a first isolation barrier 280B between the VMs 300A and 300B and a second isolation barrier 280A between the virtualization software 200A and all VMs in the computer system 2B, including the VMs 300A and 300B. The isolation barriers 280A and 280B may be referred to as “virtualization barriers” because they are implemented by the virtualization software 200A, in conjunction with the system hardware 100B. The isolation barriers 280A and 280B may also be referred to as virtualization barriers because they are established through the virtualization of hardware resources, such as the virtualization of system memory.\nIt is widely recognized that virtualization techniques can generally provide better security and more effective isolation between multiple software modules than general OSs provide. Thus, the virtualization barriers 280A and 280B of FIG. 2 can generally provide much better isolation between the multiple VMs 300A and 300B and the virtualization software 200A than the OS isolation barriers 80A and 80B of FIG. 1 provide between the multiple applications 40A and 40B and the OS 20A. This improved isolation can be attributed to a variety of factors, depending on the particular situation. For example, OSs, such as the OS 20A, generally provide an Application Programming Interface (API) for use by applications, such as the applications 40A and 40B. These APIs can be used by application software first to discover and then to exploit vulnerabilities within the OS. Virtualization software, in contrast, generally has no such interface, or at least a more limited and/or restricted interface, for use by the software within a VM. Also, virtualization software generally does not need to rely on the integrity of third party device drivers. In contrast, general OSs often allow third party device drivers to execute at a most privileged level on the system hardware, so that they can access restricted hardware resources. If such a device driver contains malicious software, the driver can compromise the integrity of the entire computer system. Also, virtualization software is generally smaller and, in some ways, less complex than a general OS, such as a Windows or Linux OS, which generally leads to a smaller number of vulnerabilities that may be exploited to compromise the integrity of the computer system. Thus, virtualization software is likely to have fewer vulnerabilities than a general OS, and, more importantly, by the very nature of virtualizing a computer system, the virtualization software is much more isolated from software within a VM than a general OS is from other software within a computer system, making it much harder to discover and exploit any possible vulnerabilities.\nAlthough computer systems that establish multiple VMs and that run different software modules within the different VMs generally provide better isolation for the software modules than do general OSs, such virtual computer systems have other limitations. First, if the software within a VM becomes corrupted by malicious software, the same problems described above relative to non-virtualized computer systems can occur within the affected VM. All software modules within the particular VM may be compromised by the malicious software. This scenario may be seen in FIG. 2. Suppose a malicious software module is able to execute within the VM 300A. If the malicious software is able to circumvent the OS isolation barriers 80A and/or 80B within that VM, then the malicious software can corrupt the OS 20B and/or the applications 40C and 40D. The virtualization barriers 280A and 280B will generally contain the malicious software within the VM 300A, however, so that it is not able to corrupt the virtualization software 200A or the software within the VM 300B. Thus, while the computer system 2B of FIG. 2 is able to limit the corruption to the VM 300A, the software within the VM 300A may be no more secure than the software within the computer system 2A of FIG. 1, depending on the circumstances. Of course, more VMs can be created within the virtual computer system 2B of FIG. 2 to further isolate the software modules within the computer system, but this increased isolation aggravates a second limitation of these virtual computer systems.\nMultiple VMs within a computer system are generally similar to multiple physical computer systems. The multiple VMs generally have, at most, only limited communication and interaction with one another. For example, in the computer system of the '774 patent, the multiple VMs in the system can communicate with each other by sending information to a server, which can send the information on to another VM in the system, except as limited by user-definable rules. In the computer system described in the Terra paper, VMs on a single physical machine can communicate with one another over virtualized standard input/output interfaces such as network interface cards, serial ports, etc. Many software modules within a computer system, however, require more effective interaction or communication with other software modules. For example, multiple software modules may need to use shared memory to facilitate fast and efficient communication of data, or a software module may need to call a subroutine of another software module.\nAnother disadvantage of a computer system in which multiple VMs are used to isolate multiple software modules is that setting up and using such a system generally takes significantly more time and effort than a computer system in which multiple software modules and a general OS run directly on physical hardware. Individuals and organizations may be reluctant to adopt a multiple VM solution or they may not use such a solution consistently enough. Anytime such a solution is not used, and multiple software modules are run directly on a physical computer system or within a single VM, the risk of corruption by a malicious software module is increased.\nWhat is needed therefore is a computer system for executing multiple software modules that provides improved security, such as is provided by the virtualized computer systems described above, while also providing more efficient and effective communication and/or interaction between multiple software modules. It would also be advantageous if such a computer system were easy to configure, use and maintain."}
{"text": "Conventional approaches for cardiac valve replacement require the cutting of a relatively large opening in the patient's sternum (“sternotomy”) or thoracic cavity (“thoracotomy”) in order to allow the surgeon to access the patient's heart. Additionally, these approaches require arrest of the patient's heart and a cardiopulmonary bypass (i.e., use of a heart-lung bypass machine to oxygenate and circulate the patient's blood). Despite their invasiveness, these surgical approaches may be reasonably safe for a first intervention. However, tissue adherences resulting from the first surgery may increase the risks (e.g., death) associated with subsequent valve replacement surgeries. See Akins et al., “Risk of Reoperative Valve Replacement for Failed Mitral and Aortic Bioprostheses”, Ann Thorac Surg 1998; 65:1545-52; and Weerasinghe et al., “First Redo Heart Valve Replacement—A 10-Year Analysis”, Circulation 1999; 99:655-658; each of which is incorporated by reference herein in its entirety.\nSynthetic valves and biological valves have been used for cardiac valve replacement with varying results. Synthetic valves rarely fail but require life-long anti-coagulant treatment to prevent blood from clotting (thrombosis) in and around the replacement valve. Such anti-coagulant treatment significantly limits patients' activities and can cause various other complications. Biological valves do not require such anti-coagulation treatment but typically fail within 10-15 years. Thus, to limit the need for and risks associated with re-operation on failed biological valves, traditionally only patients with less than about 10-15 years to live have received biological valve replacements. Patients with longer life expectancies have received synthetic valves and anti-coagulant treatment.\nAttempts have been made to develop less-invasive surgical methods for cardiac valve replacement. These surgical methods, referred to as percutaneous heart valve replacement therapies (PHVT), use a catheter to deliver a replacement valve to an implantation site using the patient's vascular system. These PHVT attempts have various shortcomings, including their inability to ensure proper positioning and stability of the replacement valve within the patient's body.\nConventional closure devices for closing access orifices are also lacking in several respects, including the looseness of their fit which can cause bleeding after surgery. These closure devices also lack a central lumen, which renders them incompatible with guide wire delivery systems. One such conventional closure device is described in Malgorzata Pawelec-Wojtalik, “Closure of left ventricle perforation with the use of muscular VSD occluder”, European Journal of Cardio-Thoracic Surgery 27 (2005) 714-716, which is incorporated by reference herein in its entirety.\nIn view of the foregoing, it would be desirable to provide improved methods, systems, and devices for cardiac valve replacement."}
{"text": "1. Field of the Invention:\nThe present invention relates generally to towing devices for transporting disabled vehicles.\n2. Description of the Prior Art:\nIn towing a disabled vehicle it is important to raise the wheels coupled to the transmission to avoid damage to the transmission. Various dollys have been proposed which are adopted to receive and capture a set of wheels for this purpose.\nSome prior art devices, such as that shown in U.S. Pat. No. 2,943,863 (Cory), require a tow truck to hoist the vehicle onto the dolly. This is a slow procedure and requires considerable room to maneuver the tow truck. Other prior art devices, U.S. Pat. No. 3,035,812 (Winteer), being examplary, require hydraulic jacks to hoist the receptacles holding the wheels of the disabled vehicle. These devices are complex, expensive to manufacture, and present the danger of hydraulic jack failure and consequent dropping of the vehicle. Furthermore, such jacks require tedious, repetitious movements of the jack handle to achieve hoisting."}
{"text": "Generally, a conventional pump-type cosmetic container has a pumping structure that discharges the contents by pumping the pumping member which is attached to the top of the container body, and a container body that holds contents, a pumping member that is connected to the top of the container body and draws up the contents by making the inside of the container vacuous and pumping the contents out, and a piston that is secured inside the container body and ascends/descends by the movement of the pump.\nThe above pump-type cosmetic containers can hold various contents in the different container depending on the purpose of the skin care, such as whitening cream and anti-wrinkle cream. Since the different cosmetics are sold in separated containers, users should purchase each of cosmetics that is contained in a separate container according to different purpose. Furthermore they find it inconvenient to carry all the separate containers in their bag, due to their volume.\nIn addition, when they go out, users should carry a lot of cosmetic containers for skin care. Even though some of cosmetics are supposed to be applied on the face in certain order, such as skin lotion followed by emulsion, they cannot provide themselves with proper skincare when they forget one or two cosmetics by carelessness. In addition, they have to check all of each container whenever they do makeup, which results in inconvenience for the users."}
{"text": "1. Field of the Invention\nThe present invention relates to an endoscope system.\n2. Description of the Related Art\nFor example, JP-A-2000-342533 discloses an endoscope system including an endoscope which has an imaging element emitting illumination light from a front end of an endoscope insertion unit and capturing an image of an observation area, and a control device which supplies the illumination light to the endoscope while being connected to the endoscope, and performs a calculation process on an imaging signal from the imaging element. In addition, recently, a technology of using a combination of a fluorescent body and laser light as the illumination light of the endoscope has been put to practical use. In the case of the endoscope system using the illumination light formed by the fluorescent body and the laser light, the fluorescent body is disposed in the front end of the endoscope insertion unit, and the laser light is supplied to the fluorescent body via a thin optical fiber cable, which is an advantage in that the endoscope insertion unit is thin.\nIn such an endoscope system, in order to accurately examine an affected portion, a chromaticity adjustment process is performed by correcting the chromaticity of a captured image to be the correct chromaticity on the basis of a predetermined chromaticity correcting table. However, there is an individual difference in the optical characteristics of each endoscope and the optical characteristics of a control device to which the endoscope is connected. For this reason, even when the uniform correction of the image signal output from the endoscope is attempted by the control device, it is difficult to perform the correction with correct chromaticity at all times. Particularly, in the case of the illumination device obtained by combining the fluorescent body with the laser light source, since the tone of the final observation image is determined by the delicate relationship of the difference in the light emission characteristics caused by the individual difference of the fluorescent body or the difference in the light emission wavelength caused by the individual difference of the illumination light source, when the endoscope connected to the control device is exchanged with another one, the tone of the captured image may be different for each endoscope.\nFIG. 12 shows a distribution of a variation in the illumination light of the endoscope caused by the individual differences of chromaticities (X, Y) of the light sources. When the laser light source is used, since the laser light emitting element has a variation in the light emission wavelength of about ±5 nm, a variation in the chromaticity caused by the individual differences of the light sources (for example, excitation light sources of 440 nm and 445 nm) is about 0.02 in the chromaticities X and Y. When a variation in the characteristics of the illumination light is small, the variation can be adjusted by a white balance function provided in a control device. However, when a variation caused by the individual differences is large, and is out of an adjustable range, the correct tone cannot be realized even when the white balance adjustment is performed. For this reason, the chromaticity of the illumination light determines the color quality of the observation image."}
{"text": "The present invention relates to a sense amplifier, and more particularly to a sense amplifier for a multilevel cell memory.\nMemory systems include sense amplifiers for detecting content stored in memory cells. In digital multilevel memories, reference cells store reference values for comparison with detected content in memory cells. The speed of and the loading on the sense amplifier impacts its performance."}
{"text": "The role of a dishwashing composition is twofold: to clean soiled dishware and to leave it shiny. Typically when water dries from surfaces water-marks, smears and/or spots are left behind. These water-marks may be due to the evaporation of water from the surface leaving behind deposits of minerals which were present as dissolved solids in the water, for example calcium, magnesium and sodium ions and salts thereof or may be deposits of water-carried soils, or even remnants from the cleaning product, for example soap scum. During the course of this work, it has been observed that this problem can be often exacerbated by some cleaning compositions which modify the surface of the dishware during the automatic dishwashing process in such a way that after rinsing, water forms discrete droplets or beads of water remain on the surface instead of draining off. These droplets or beads dry to leave noticeable spots or marks known as water-marks. This problem is particularly apparent on ceramic, stainless steel, plastic, glass and painted surfaces.\nWhen items are washed in an automatic dishwasher it is important to the consumer that those items come out at the end of the cycle as clean as possible. An item is clean to the consumer if there are no visible pieces of soil or hard water deposits. This means no films, spots, grit or residues of anything. It is also as important that when the consumer touches the items they do not feel anything other than the clean surface of item. When an item is not clean, a consumer may feel either a rough and gritty surface or they may feel a greasy surface.\nA rough and gritty surface can be produced when hard water deposits, i.e. calcium carbonate and other salts deposit on the item. This can be accentuated if there are food soils mixed in with the deposits. A greasy surface can be produced if excess greases and fat soils from the wash liquor of the dishwashing process have been deposited onto the item. Whether rough and gritty or greasy, these feelings are unpleasant for the consumer and indicate that the items which they have washed are not clean.\nThe object of the present invention is to provide a dishwashing composition that leaves the washed dishware clean with reduced incidence or spots."}
{"text": "1. Field of the Invention\nThe present invention relates to a scanning head having a structure suitable for a printer, scanner, copier or other image input and/or output apparatus, and a printer having such a scanning head.\n2. Description of the Related Art\nPage printers have been vigorously developed in recent years, because they can print on ordinary paper as well as specific paper. A page printer uses a laser scanning head composed of a laser diode and a polygon lens. In a laser scanning head, a laser emitting point is moved by a polygon lens, and printing at high-speed is difficult.\nFor high-speed printing, an LED scanning head using two or more LEDs has been developed. Two or more LEDs are aligned in an LED scanning head. These LEDs simultaneously emit light of different intensity, thereby scanning a photoconductor. As high picture quality is demanded, very high accuracy is demanded for packaging of LEDS with high density. This causes a problem of increased number of components.\nTo solve the above problems, Jpn. Pat. Appln. KOKAI Publication No. 9-226172 proposed a scanning head using an organic electroluminescent element as an LED.\nHowever, at present, an organic electroluminescent element has a problem in luminous intensity and life. Namely, a light-emitting element requires sufficient amount of light to expose a photoconductor, and if the luminous intensity per dot of an organic electroluminescent element is weak, the exposing time per a dot must be set long. To set the exposing time long, the printing speed must be delayed. Conversely, if the luminous intensity per a dot of an organic electroluminescent element is increased, the exposing time per dot is reduced and the printing time is reduced, but the life of an organic electroluminescent element is reduced.\nThe luminous flux of an LED such as an organic electroluminescent element spreads from a light-emitting point, and it is preferable to provide an optical system between LED and photoconductor, which gives directivity to incident light from a dot of an LED to be emitted only to a specified part of a photoconductor. However, the efficiency of using such an optical system depends on an angle of taking in light (angular aperture), and the efficiency of using light is not high in a source like an LED which causes a light diffusion."}
{"text": "This new hybrid tea rose cultivar originated as a seedling resulting from a crossing of an unnamed seed parent with an unnamed pollen parent. The seed parent was derived as follows: (Sterling Silver.times.Intermezzo).times.(Sterling Silver.times.Simone). Sterling Silver is covered by U.S. Plant Pat. No. 1,433, Intermezzo by U.S. Plant Pat. No. 2,430, and Simone by U.S. Plant Pat. No. 1,847. The pollen parent was derived as followed: (Music Maker).times.(Mainzer Fastnacht.times.Tom Brown). The new variety was discovered by me in 1975, at Iowa State University Horticulture Greenhouses, Ames, Iowa, in the course of breeding efforts to produce everblooming strains of hybrid tea roses over a wide range of flower colors and a capability of withstanding winter cold. Such breeding efforts began in about 1949 and were carried out by me at the University in Ames, Iowa.\nFollowing its discovery, this new plant has been asexually reproduced under my direction at Ames, Iowa and Tyler, Tex. by cuttings and by bud-grafting. Such propagation through successive generations has shown that the distinctive characteristics of the new variety hold true from generation to generation and appear to be firmly fixed."}
{"text": "Modern phone lines are being upgraded from analog circuits in increasing numbers to voice over Internet protocol (VoIP). VoIP is a technology that allows a caller to make voice calls using a broadband Internet connection instead of a conventional (or analog) phone line. VoIP services convert voice audio into digital signals that travel over a network, which may often include segments considered part of the Internet. When calling a regular phone number from a VoIP number, the signal is converted to a regular telephone signal before it reaches the destination. VoIP calls can be made directly from a computer, a special VoIP phone, or a traditional phone connected to an adapter. In addition, wireless “hot spots” in locations such as airports, parks, and cafes with connection to the Internet may enable use of VoIP service wirelessly.\nWith its portability, cost savings, and the promise of enhanced functionality, VoIP telephony is becoming the de facto choice in many applications. VoIP has a range of advantages. It is a flexible communication system. For example, an end user assigned a VoIP line associated with a New York telephone number can move to a different state, e.g., Colorado, and still receive calls to the same New York telephone number associated with the VoIP line despite the end user physically relocating to Colorado. As such, the VoIP telephone number can travel with the end user such that the end user can be accessible via a single VoIP telephone number at multiple locations.\nHowever, the virtual nature and portability of VoIP can cause various issues in the context of call routing, billing, and regulations compliance. Specifically, conventional systems route VoIP calls through telecommunication networks based upon the VoIP telephone number itself, i.e., the ten digit phone number (in the United States) and as such are ill equipped to route and process VoIP calls when end users operate VoIP lines outside of the geographical areas associated with the respective VoIP telephone numbers.\nAccordingly, a need exists for improved routing devices and systems to address the portability of VoIP, among other improvements. It is with these observations in mind, among others, that various aspects of the present disclosure were conceived and developed."}
{"text": "The present invention is directed to apparatus for monitoring engine exhaust, comprising a chamber for receiving such exhaust, an inlet of the chamber adapted to be connected to the engine exhaust and an outlet from the chamber to enable exhaust to escape therefrom, so that exhaust from the engine can flow through the chamber continuously, a source of electromagnetic radiation which is of a wavelength within a range comprising the infrared, visible and ultraviolet wavelengths, detector means which detects the amount of such radiation it receives, the source and the detector being so arranged that radiation emitted by the source is received by the detector after the radiation has passed through the chamber, and electrical means connected to the detector to provide an electrical signal which is indicative of the attenuation of the radiation as a result of the presence of particles within the exhaust in the chamber, in which the detector means is such as to be able to provide respective different measures of radiation which it receives for at least two different wavelengths of such radiation, and the electrical means are such as to provide respective electrical signals which are indicative of the attenuation of the radiation, as a result of the presence of particles within the exhaust in the enclosure, for the said at least two different wavelengths.\nOne shortcoming of such apparatus is that it takes no account of the effects of nitrogen dioxide present in the exhaust. This is a brown gas which absorbs radiation in much of the visible and UV range. Its concentration may fluctuate, and its opacity as a function of wavelength is not simple. Furthermore, oxidation catalysts now fitted to many vehicles increases the amount of nitrogen dioxide present in the engine exhaust as much as tenfold. At the same time, engines are becoming more efficient, so that soot concentration levels have become low. These things very much increase the adverse effect of nitrogen dioxide on the apparatus."}
{"text": "1. Field of the Invention\nThe present invention relates to an inkjet printing apparatus and an inkjet printing method which use color inks containing colorants and an image quality improvement liquid, and more particularly to a technology for reducing gloss unevenness in printed images.\n2. Description of the Related Art\nThere has been growing calls in recent years for the inkjet printing to have a capability to print high quality images on a variety of print mediums. Among the print mediums suited for high quality images, there is coated paper. The coated paper has an ink receiving layer formed on a substrate such as quality paper and film. There are various kinds of coated paper with varying degrees of texture, from glossy paper with a mirror surface to matte paper with a glare-free finish.\nFor these coated paper, there is a wide range of demands in terms of glossiness of printed images. One such demand is that the printed image be uniform in glossiness over the entire print medium. To meet this demand Japanese Patent No. 4003760 discloses a method that, in an inkjet printing apparatus using color inks and an image quality improvement liquid, alleviates gloss unevenness by adjusting the amount of image quality improvement liquid applied according to the volume of color inks used for printing.\nGenerally, in areas on glossy paper applied with a small volume of color inks, the level of gloss, which will be described later, is low compared with areas applied with a greater amount of inks. So, Japanese Patent No. 4003760 minimizes the gloss unevenness within the same image by applying a greater amount of image quality improvement liquid to the areas printed with a small volume of inks than to those areas printed with a larger volume of inks to enhance the level of gloss in the areas printed with a small ink volume.\nHowever, as disclosed in Japanese Patent No. 4003760, with the method of making only the gloss level uniform by adjusting the amount of image quality improvement liquid, the uniformity of glossiness in the same image may not be able to be enhanced enough. This is considered due to the fact that the glossiness in an image is affected by not only the uniformity of gloss level but the uniformity of image clarity and that the image clarity and the gloss level change according to the gradation value of the printed areas.\nFIG. 1 illustrates how the image clarity and the gloss level vary according to the gradation value. In FIG. 1, “medium” represents a target range of each of the image clarity and the gloss level; “high” represents a range higher than the target range; and “low” represents a range lower than the target range. As shown in FIG. 1, there is a tendency that, when compared with the target range, highlight areas are high in image clarity and low in gloss level, halftone areas are high in both image clarity and gloss level, and shadow areas (high density areas) are medium in image clarity and high in gloss level. This shows that the image clarity as well as the gloss level tends to vary according to the gradation value, which means that the user can recognize gloss unevenness when the image clarity uniformity is low even if the gloss level is uniform. The gloss unevenness in an image becomes particularly distinctive when the gloss level and the image clarity greatly vary between highlight areas and shadow areas."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for preparing a low-odor, and storage-stable, monomer-containing polyisocyanurate based on isophorone diisocyanate.\n2. Discussion of the Background\nAs polyisocyanate adducts, polyisocyanurates are valuable components in the production of high-grade coatings having good mechanical properties and also good light stability and weather resistance. Polyisocyanurates based on isophorone diisocyanate (IPDI) are also used as raw materials for PU-based elastomer applications. It may be desirable for the IPDI-based polyisocyanurate, also referred to as IPDI trimer, to be used in a form which still includes monomer.\nPolyisocyanurates are basically obtained by catalytic trimerization of appropriate isocyanates. Examples of appropriate isocyanates are aromatic, cycloaliphatic, and aliphatic diisocyanates and higher polyisocyanates. Examples of suitable catalysts include tertiary amines (U.S. Pat. No. 3,996,223), alkali metal salts of carboxylic acids (CA 2113890; EP 56159), quaternary ammonium salts (EP 798299; EP 524501; U.S. Pat. Nos. 4,186,255; 5,258,482; 4,503,226; 5,221,743), aminosilanes (EP 197864; U.S. Pat. No. 4,697,014), and quaternary hydroxyalkylammonium salts (EP 17998; U.S. Pat. No. 4,324,879). Depending on the catalyst it is also possible to use various cocatalysts, for example, OH-functionalized compounds or Mannich bases composed of secondary amines and aldehydes and/or ketones.\nFor trimerization, the polyisocyanates are reacted in the presence of the catalyst, where appropriate using solvents and/or auxiliaries, until the desired conversion has been achieved. In this context, one also speaks of partial trimerization, since the target conversion is generally well below 100%. The reaction is then terminated by deactivating the catalyst. This is done by adding a catalyst inhibitor such as, for example, p-toluenesulfonic acid, hydrogen chloride or dibutyl phosphate, and results unavoidably in a possibly unwanted contamination of the polyisocyanate containing isocyanurate groups that are formed. Of particular advantage in respect of the trimerization of isocyanates on the industrial scale is the use of quaternary hydroxyalkylammonium carboxylates as oligomerization catalysts. This type of catalyst is thermally labile and is amenable to targeted thermal deactivation, so that it is unnecessary to stop the trimerization by adding potentially quality-reducing inhibitors when the desired conversion has been reached.\nMonomer-containing IPDI trimer, which is suitable, for example, for the abovementioned PU injection applications, has an NCO content of 28-32% by weight. The polyisocyanurate is prepared by partial trimerization of IPDI in the presence of one or more appropriate catalysts. Afterward, the catalyst must either be removed fully from the reaction solutionxe2x80x94this can be done by short-path or thin-film distillationxe2x80x94or deactivated, because the trimer lacks storage stability in the presence of active catalyst residues. Where the NCO content of the IPDI polyisocyanurate obtained is below the desired level, it can be adjusted as desired without problems by diluting the solution with monomeric IPDI.\nAlkali metal salts of carboxylic acids are not very suitable as catalysts for preparing monomer-containing IPDI trimer, because they are difficult, if not impossible, to remove from the reaction product. With regard to the amine catalysts available, it has been found that the resulting IPDI trimer solutions are fundamentally hampered by a clearly perceptible odor, which is sufficiently pronounced to remain detectable in the end application and to manifest itself unpleasantly. In order to eliminate the odor nuisance, it is technical practice to free the reaction solution, following partial trimerization and catalyst deactivation, from excess IPDI, from odoriferous components, and, where appropriate, from unwanted catalyst inhibitors. This is generally done by means of short-path or thin-film distillation. Subsequently, the monomer-freed solid resin is converted into the desired, low-odor and monomer-containing IPDI polyisocyanurate by adding fresh IPDI.\nThe sequence of partial trimerization/deactivation, demonomerization/purification, and, finally, dissolution of the solid resin in the monomer is very cumbersome. The step of separating off the monomer in the existing process consumes a lot of time, adds a lot of costs, and, moreover, represents a capacity-limiting bottleneck.\nIt is therefore an object of the present invention to provide a more economical process for preparing a low-odor and storage-stable monomer-containing polyisocyanurate based on isophorone diisocyanate, in which demonomerization is avoided. It has surprisingly been found that this step can in fact be omitted and, moreover, that the use of possibly quality-reducing inhibitors can be avoided if the trimerization of IPDI is conducted in the presence of specific catalysts.\nAccordingly, the present invention provides for a process for preparing a low-odor, storage-stable monomer-containing polyisocyanurate, comprising:\npartially trimerizing isophorone diisocyanate at a temperature of 0-160xc2x0 C. for 3 minutes to 3 hours in the presence of 0.05-2% by weight of a catalyst, based on the weight of the isophorone diisocyanate;\nwherein said catalyst is represented by the following formula\n[Rxe2x80x94NX3]m⊕mYxe2x8ax96\nwherein\nYxe2x8ax96 is a carboxylic acid anion having 4-8 carbons;\nR is a xcex2-hydroxyalkyl group having 2-6 carbons;\nX is an alkylene group having 2-3 carbons; and\nm is a number from 1.0 to 2.0;\nwherein the three radicals X form a ring with the quaternary nitrogen by way of a common nitrogen atom; and\nwherein said process proceeds at a temperature of from 0 to 160xc2x0 C.\nThe present invention provides a process for preparing a low-odor, storage-stable monomer-containing polyisocyanurate from isophorone diisocyanate. The process comprises conducting a partial trimerization of isophorone diisocyanate over a period of 3 minutes to 3 hours in the presence of 0.05-2% by weight of a catalyst, based on the weight of the diisocyanate, the catalyst having the general formula\nxe2x80x83[Rxe2x80x94NX3]m⊕mYxe2x8ax96\nwherein Yxe2x8ax96 is a carboxylic acid anion of 4-8 carbons,\nR is a xcex2-hydroxyalkyl group of 2-6 carbons,\nX is an alkylene group of 2-3 carbons, and\nm is a number from 1.0 to 2.0. The three radicals X form a ring with the quaternary nitrogen by way of a common nitrogen atom, which may be partly xcex2-hydroxyalkylated, said ring possibly having an OH group positioned xcex1, xcex2 or xcex3 to the nitrogen. The process proceeds at a temperature of 0-160xc2x0 C., preferably 40-120xc2x0 C., and most preferably 55-95xc2x0 C. In this it is possible to dispense with separating off the monomer and chemically deactivating the trimerization catalyst. The reaction temperature includes all values and subvalues therebetween, especially including 10, 20, 30, 40, 50, 60, 70, 80, 90, 100, 110, 120, 130, 140 and 150xc2x0 C. The reaction time of the partial trimerization includes all values and subvalues therebetween, especially including 5, 10, 20, 30, 40, 50, 60, 70, 80, 90, 100, 110, 120, 130, 140, 150, 160 and 170 min. The amount of catalyst includes all values and subvalues therebetween, especially including 0.1, 0.2, 0.3, 0.4, 0.5, 0.6, 0.7, 1, 1.1, 1.2, 1.3, 1.4, 1.5, 1.6, 1.7, 1.8 and 1.9%.\nIn principle, isocyanates appropriate for trimerization can be prepared by different kinds of processes (Annalen der Chemie 562 (1949) 75 ff.). Particularly well established in the industry is their preparation by phosgenating organic polyamines to the corresponding polycarbamoyl chlorides and cleaving these chlorides thermally into organic polyisocyanates and hydrogen chloride. Alternatively, organic polyisocyanates can also be prepared without using phosgene, i.e., by phosgene-free processes. According to EP-A-126 299 (U.S. Pat. No. 4,596,678), EP-A-126 300 (U.S. Pat. No. 4,596,679), and EP-A-355 443 (U.S. Pat. No. 5,087,739), for example, (cyclo)aliphatic diisocyanates, such as 1-isocyanato-3-isocyanatomethyl-3,5,5-trimethylcyclohexane (isophorone diisocyanate, IPDI), can be made available by reacting the parent (cyclo)aliphatic diamines with urea and alcohols to give (cyclo)aliphatic biscarbamic esters, and thermally cleaving these esters into the corresponding diisocyanates and alcohols.\nFor the process of the present invention for preparing low-odor, storage-stable monomer-containing polyisocyanurates based on isophorone diisocyanate it is unimportant by which synthetic pathway the IPDI used has been prepared. It may be noted, however, that one of the factors affecting the amount of catalyst required for achieving a desired NCO content is the quality of the raw material. The NCO is 22-34% by weight, especially including 23, 24, 25, 26, 27, 28, 29, 30, 31, 32 and 33% by weight. From experience, increasing levels of hydrolyzable chlorine compounds in the IPDI necessitate an increase in the amount of catalyst. It appears that the hydrolyzable chlorine tends to have an inhibiting effect on the catalyst.\nThe preparation of the low-odor, storage-stable monomer-containing polyisocyanurates based on isophorone diisocyanate by partial trimerization, in accordance with the invention, may take place continuously (tube reactor or tank cascade) or else may be conducted batchwise. The catalyst is used at a low concentration of between 0.05 and 2% by weight. The precise amount can be determined experimentally with ease and is dependent on the catalyst, on the target conversion, on the quality of the IPDI used, and on the procedure.\nThe partial trimerization may be conducted within a period of from 3 minutes to 3 hours. The product contains monomeric IPDI, trimeric IPDI isocyanurate, and higher oligomers with an isocyanurate structure. As a secondary component, minor amounts of compounds with a uretdione structure may also be found. Compounds of this kind are described in the literature.\nIn accordance with the present invention the catalyst is used in an amount of 0.05-2% by weight, based on the weight of the isophorone diisocyanate used. The catalyst is easy to obtain by reacting the three building blocksxe2x80x94tricyclic diamine, carboxylic acid, and oxirane. Where appropriate, it may be prepared in the presence of a solvent. It is common to use alcohols of low molecular mass such as methanol or ethylene glycol. One example of a suitable tricyclic diamine is diazabicyclo[2.2.2]octane. Examples of suitable carboxylic acids are acetic acid, hexanoic acid, and 2-ethylhexanoic acid. Possible options for the oxirane component include, for example, propylene oxide, butylene oxide or 1,2-epoxyhexane. The molar ratios of the three building blocks for preparing the catalysts used in accordance with the invention are variable. Depending on the proportions of the building blocks, catalysts having at least one quaternary nitrogen atom in the molecule are obtained. As the examples show, catalysts possessing in some cases more than two quaternary nitrogen atoms are also used.\nThe process of the present invention is conducted at temperatures between 0xc2x0 C. and 160xc2x0 C., preferably between 40xc2x0 C. and 120xc2x0 C., and, with particular preference, between 55 and 95xc2x0 C.\nIn accordance with the present invention, the partial trimerization of the isophorone diisocyanate is conducted either batchwise or continuously. The batch process is preferred.\nIn the case of the batch process, operation takes place in a stirred reactor. The isophorone diisocyanate is introduced as initial charge. The catalyst is metered in after the IPDI has reached the temperature required for reaction: 0-140xc2x0 C., preferably 55-90xc2x0 C., particularly preferably 65-75xc2x0 C. The temperature includes all values and subvalues therebetween, especially including 10, 20, 30, 40, 50, 60, 70, 80, 90, 100, 110, 120 and 130xc2x0 C. The trimerization is exothermic. The catalyst is first of all metered in in an amount such that a marked rise in the temperature of the reaction mixture by 5-15xc2x0 C. occurs. The rise in temperature includes all values and subvalues therebetween, especially including 6, 7, 8, 9, 10, 11, 12, 13 and 14xc2x0 C. Since the catalyst is deactivated in the course of the reaction, the temperature of the reaction mixture drops again during the reaction, and a further addition of catalyst may be made. This procedure is repeated until the desired conversion has been reached. The sequence of catalyst deactivation and reinitiation of the trimerization by addition of further portions of catalyst means that it is readily possible at any time to monitor the process in respect to both, the conversion and the temperature profile of the reaction.\nThe catalyst may be used in neat form. For more precise metering and optimum mixing of the catalyst, however, it is advantageous to dissolve the catalyst in an appropriate solvent. Appropriate solvents include in principle all those possessing good solvency for the catalyst: examples include water, alcohols of low molecular mass such as methanol or ethylene glycol, or organic acids of low molecular mass such as acetic or hexanoic acid, for example.\nThe continuous trimerization may be conducted in a tank cascade. A combination of tank cascade and tube reactor is also possible.\nIn order to restrict the required amount of catalyst relative to the desired conversion, the temperature profile of the process of the invention should as far as possible be configured so that the reaction solution does not exceed a temperature of 95xc2x0 C., preferably 85xc2x0 C., more preferably 75xc2x0 C. and most preferably 65xc2x0 C.\nThe low-odor, storage-stable monomer-containing polyisocyanurates based on isophorone diisocyanate that are prepared in accordance with the present invention are useful intermediates for polyurethane coatings, polyurethane and polyurea moldings, as produced, for example, by the RIM (reaction injection molding) process, for polyurethane injection applications, or else for polyurethane-based automobile window sealing. They may be used in a form in which they have been blocked with blocking agents. Examples of suitable blocking agents in this case include lactams, such as xcex5-caprolactam, oximes, such as methyl ethyl ketoxime or butanone oxime; triazoles, such as 1H-1,2,4-triazole; readily enolyzable compounds, such as acetoacetates or acetylacetone; or else malonic acid derivatives, such as malonic diesters. The blocking agents can be used alone or in combination.\nHaving generally described this invention, a further understanding can be obtained by reference to certain specific examples which are provided herein for purposes of illustration only, and are not intended to be limiting unless otherwise specified."}
{"text": "It is known to provide sliding doors to control access to rooms of one type or another including shower recesses, however with current mounting assemblies and hardware it is difficult to maintain the mounting assemblies in a good working order and in good appearance, due to build up of grime, dirt, dust whilst still affording a degree of protection to the mounting assemblies.\nThe subject application seeks to ameliorate one or more of the abovementioned disadvantages."}
{"text": "The present day digital camera includes many functions. For example, there is known a digital camera which has a character photographing mode in addition to a normal photographing mode. The normal photographing mode is the mode in which ordinary snapshots, for example, photographs of scenery are taken. In the character photographing mode, the digital camera subjects data to image processing so as to improve the appearance of the characters.\nThere is a digital still camera disclosed in Japanese patent application laid open No. HEI 08-125870. This digital still camera has a structure as follows. The digital still camera has the normal photographing mode and the document photographing mode. Any one of these two modes can be selected. When the normal photographing mode is selected, image data is compressed by a natural picture compression unit. Whereas image data is compressed by a document compression unit (a compression method appropriate for an image with a small number of gray levels in achromatic color) when the document photographing mode is selected. The compressed image data is then stored in an externally provided memory. Therefore, the image data can efficiently be compressed irrespective of whether the data belongs to a picture or belongs to a photograph of a document without degradation in image quality.\nHowever, the conventional type of digital camera has a problem that the data of the captured image can not effectively be used. More specifically, the digital camera having the character photographing mode can not transmit the captured image as through a facsimile. In order to transmit the image as a facsimile, following process is carried out. First, the image is transmitted to a computer where the image is subjected to processing such as expansion, image conversion, compression, or the like and then only the image can be transmitted as a facsimile.\nFurther, in the digital still camera described above, the compression method is uniquely determined based on the photographing mode. For example, when a document (for example, an encyclopedia) which contains picture of scenery is photographed, and if the picture is photographed in the document mode, the image data is compressed by the compression method suitable for an image with a small number of gray levels in achromatic color. Resultantly, the picture of the scenery can not be reproduced properly. Therefore, the user is required to photograph the same document in the normal photographing mode. Thus, the user has to choose the photographing modes based on the purpose of using the document."}
{"text": "Agrochemicals, water-treatment chemicals, detergents and the like are commonly supplied in unit-dose pouches where the active chemical is contained within a water-soluble (usually polyvinyl alcohol) film. Increasingly regulatory authorities are requiring that each of these unit-dose soluble packages are individually labelled. Thus, for example, it is proposed that in the EU that any detergent unit-dose package containing more than 25 ml of an irritant or corrosive chemical will need to be individually labelled as such.\nPrinting on the individual pouches may be done either on the web of film before the pouch has been made, when the image can be printed on either side of the film so that it ends up inside or outside the finished pouch, or after the pouch has been made when the image is on the outside of the pouch.\nDetergent pouches can be produced either on rotary drum machines such as those produced by Cloud Packaging Equipment, Eme Engel Machinefabriek en Engineering B.V. or Green Sustainable Packaging Inc. Pouches can also be produced on a flat bed machine such as those produced by Harro Höfliger Verpackungsmaschinen GmbH. On rotary drum machines ink-jet printers can be installed on either one of the webs feeding the machine or placed to print on the formed pouches.\nThe colorants for use in printing on these water-soluble materials must meet many demanding criteria. Thus, they must display excellent adhesion to water-soluble materials (such as polyvinyl alcohol), they must give clear images with good scuff and scratch resistance and crucially they must be able to readily disperse/dissolve in water without leaving any residue or ink skin. It has been found that pigments give excellent performance.\nSince the colorant must dissipate when the water-soluble pouch dissolves the pigments (which are in their native state insoluble) must be in a form able to readily disperse.\nThere are a number of different treatments which enable pigments to be self-dispersible. However, this technology is usually aimed at providing a pigment which can form a stable dispersion in a defined liquid medium. Dispersing a self-dispersible pigment in a defined liquid medium, such as an ink or paint, typically requires careful processing and the selection of solvents and other additives able to help disperse and stabilise the pigment.\nThis is a very different challenge to providing a pigment able to completely and rapidly self-disperse at room temperature in an aqueous environment and existing commercially available self-dispersible pigments are not able to satisfactorily disperse without leaving an undesirable residue."}
{"text": "This invention relates to fiber optic devices and, more particularly, to devices for monitoring, switching, attenuating or distributing optical signals in fibers.\nOptical transmission systems are being developed to handle large amounts of communications traffic. One type of optical transmission system uses optical fibers as the transmitting media. Optical fibers with sufficiently low loss for practical long distance optical communication systems have been developed only recently.\nAny communication system requires both transmitting and receiving apparatus. It is desirable that the transmitting source have a constant power output that is independent of temperature changes and component aging. One source of optical power is a laser which changes output power in response to changes of both temperature and aging. Output power of the laser can be stabilized by monitoring the light actually coupled into the fiber and by using such light to develop a signal for feedback control of the laser. To effectively monitor such light in the fiber for controlling the laser, a tap should monitor the transmitted light without altering the characteristic of the transmitted power. All transmitted modes should be monitored equally because any substantial preference of one mode or modes over another mode or modes is unsuitable for accurate feedback control of the laser. Available fiber optic taps substantially prefer one optical mode or modes over another mode or modes.\nFiber optic systems need optical tap devices for economically monitoring and controlling the optical power. A wide variety of devices are being developed to fill these needs. Development of a device may result in a specific and perhaps optimum component for each particular application. Such a specific development can be costly. As an alternative a designer may desire to convert the optical signal to an electrical signal and process the latter electronically taking advantage of well developed and versatile electronic technology. This electronic approach introduces noise into the system as well as requiring a source of external power.\nProblems created by developing specific components for each application and by converting optical signals to electronic signals for processing are avoided by the use of an optical fiber access port in a tap for monitoring and controlling optical power and for building many optical control devices without substantial further development effort regardless of changes introduced by new fiber materials and structures.\nA practical optical communications system requires at least a few different optical devices for establishing desirable interconnection arrangements. Optical switches can be used so that part of the system can be rearranged conveniently either for testing or for providing a variety of communications needs. Additionally, a variable optical attenuator can be used for adjusting the magnitude of power being transmitted in any part of the system. Such attenuators are helpful in setting levels for testing and for establishing some branching arrangements. Another device which facilitates the design of optical communication systems is a multibranch distributor. Such a distributor divides incident light, transmitted along one fiber, into portions for further transmission along two or more other fibers.\nAs in the case of the optical fiber tap, the variable attenuator and the multiple branch distributor should perform their designated functions without substantial preference of some mode or modes over another mode or modes. Accurate transmission of information often requires faithful communication of the input signal characteristic throughout the system.\nAvailable fiber optic devices are relatively complicated and expensive devices whereas it is desirable for reasons of cost and ease of operation to have simple, inexpensive fiber optic devices. Such simple, inexpensive devices enhance the practicality of optical communications systems arrangements.\nThese problems of complicated design, high cost and mode selectivity of the existing fiber optical devices are solved by a group of fiber optic devices, each including an access port.\nIt is an object to provide a family of simple fiber optic devices for coupling optical power between fibers.\nIt is a further object to provide fiber optic devices which are substantially insensitive to optical modes carried by a fiber.\nIt is another object to provide a basic unit for constructing different optical control devices."}
{"text": "1. Field of the Invention\nThe present invention relates to compressors, pumps, and engines. More particularly, the present invention relates to a pumping apparatus that includes two housing or rotor sections that engage a spherical bearing that enables each housing section to rotate together but about different axes of rotation. These axes intersect to form an obtuse angle. Valved pistons on the housing sections pump fluid as the housing sections are rotated.\n2. General Background of the Invention\nThe three predominate forms of pumping, driving and compressing that are available on the market at the time of this document are reciprocating, mechanical screw and rotary and centrifugal.\nThe following patent documents are incorporated herein by reference:\nU.S. Pat. Nos.: 3,945,766; 4,277,228; 4,858,480; 5,249,512; 5,647,729; 6,352,418; 6,368,072; JP 02305381A and US2001/0014288.\nU.S. Published patent application Ser. No. US2001/0014288 discloses a pump with a back and forth piston motion (see FIG. 12)."}
{"text": "The conventional surge protection technique applicable to electric products can be understood by referring to the Taiwan patent publication No. 518808 “Power Socket Module with a Surge Protection Circuit”. It is characterized in that a surge protection circuit is installed in a power cord socket and electrically connected between two alternating current (AC) electrodes of the power cord socket in parallel. The surge protection circuit is used for eliminating the electric pulse between the two AC electrodes. The surge protection circuit is substantially described in the patent description and claims 4 and 5. A varistor and a fuse are connected between the two AC electrodes in serial. Wherein, when the voltage of the electric pulse exceeds a predetermined value, the varistor will conduct to eliminate the electric pulse. If the passing current of the varistor exceeds a predetermined value, the current fuse will fuse, or a light emitting diode will be off to indicate the damage of the surge protection circuit.\nAnother conventional surge protection device is described in the Taiwan patent publication No. M297037 “Lightening Strike Protector with a Protection Device”. It is characterized in the lightening strike protection element and the point discharging module mentioned in the patent description and claim 1. When the pulse voltage is higher than the preset voltage of the point discharging module, the lightening protection element is activated.\nTherefore, the principle of the above-mentioned two conventional surge protection techniques is substantially based on the existence of a core element: the varistor (or the lightening strike protection element). The varistor is capable of changing its resistance according to the voltage bridging over two ends of the element. According to such a principle, the resistance of the varistor is very high in a normal voltage. However, when the voltage increases to a certain value, the resistance will drop down rapidly in an extremely short time. In a general case, the response time is about several nano-second (ns). Therefore, when the varistor is applied to the surge protection, it will be connected to the front end of the electric product in parallel. Consequently, when the normal operation voltage is supplied to the electric product, the varistor will be in an open-circuit state. However, when the voltage rises abnormally, the varirstor will be in a short-circuit state to abuse the energy in front of the electric product, so that the electric pulse cannot influence the electric product."}
{"text": "(1) Field of the Invention\nThe present invention relates to a Ni-based alloy having better high-temperature strength and more excellent ductility as an ordinary casting material, and a turbine rotor and stator blade for a gas turbine using the Ni-based alloy.\n(2) Description of Related Art\nIn recent years, due to increase in environmental awareness such as saving of fossil fuels, reduction in the emission amount of carbon dioxide, and prevention of global warming, thermal efficiency of an internal combustion engine has been improved. It is known that thermal engines such as gas turbines and jet engines can enhance thermal efficiency most effectively by being operated with the high temperature side of the Carnot's cycle set at a higher temperature. With elevation of the temperature of the turbine inlet, importance of improvement and development of the material for use in high-temperature components of gas turbines, that is, combustors and turbine rotor or stator blades has been increased.\nAs for the material, in order to take measures to the elevation of the temperature, a Ni-based heat-resistant alloy which is more excellent in high-temperature strength is applied to rotor blades, and many Ni-based alloys are used for the rotor blades at present. However, a Co-based alloy which is excellent in corrosion resistance and weldability is used for the heat-resistant material used in gas turbine stators.\nDue to elevation of combustion temperature (turbine inlet temperature) along with improvement in efficiency of recent years, use of a Ni-based alloy which is more excellent in high-temperature strength than a Co-based alloy has been considered. In order to enhance the strength of a Ni-based alloy, large amounts of W, Mo, Ta, Co and the like which are solid solution strengthening elements are added. In addition, by adding Al, Ti and the like to utilize precipitation strength of a γ′Ni3(Al, Ti) phase which is a reinforcement phase, the Ni-based alloy has excellent high-temperature strength and thermal fatigue characteristics.\nHowever, since a Ni-based alloy is inferior to a Co-based alloy in workability such as weldability and ductility, it is difficult at present to apply a Ni-based alloy to gas turbine stator blades.\nFor example, an alloy which has a large precipitation amount of a γ′ phase and is excellent in strength characteristics is excellent in the high-temperature strength and creep characteristics, but is significantly reduced in ductility. In contrast with this, an alloy (U.S. Pat. No. 3,720,509) which is improved in workability such as ductility by reducing the precipitation amount of a γ′ phase is required to improve in strength corresponding to elevation of temperature in order to be applied to a rotor blade."}
{"text": "1. The prior art\nThe U.S. Pat. Nos. 3,387,880, 3,820,534 and 3,918,440 and the corresponding U.K. Pat. No. 1,311,519, 1,393,701 and 1,393,703 describe arrangements for the healing of compound bone fractures, comminuted fractures in conjunction with damage to soft parts of the body, pathological fractures, spontaneous fractures and for overcoming difficulties in healing due to pseudoarthroses, osteomyelitides and the like, in the case of which a magnetic field is induced to produce bioelectric potentials and currents in the defective tissue and as a result to accelerate regeneration and to enable healing to take place in cases which are resistant to therapy.\nFor this treatment a coil is implanted in the vicinity of the fracture which represents the secondary inductance of a transformer and it is associated with a osteosynthetic part of a conventional type as for example a support plate, a pin driven into the marrow of the bone or even a joint prosthesis, or a coil can be incorporated in such an osteosynthetic part of a joint prosthesis. The inducing magnetic field is excited in a coil surrounding the diseased tissue region. The coil is supplied from a low frequency generator with current and it constitutes the primary inductance of the transformer. The low frequency magnetic field of the coil pervades the whole tissue region, which comprises the bone defect and induces an electrical voltage in the implanted coil, the voltage being caused to act in the tissue via electrodes in the form of bone screws, wires or via metallic foil arranged at the surface of an implantate in an insulated manner.\nExperiments on animals and numerous clinical applications of this technique show that in the tissue pervaded by the magnetic field the bone formation prompted by the low frequency alternating electric potentials begins in the immediate neighborhood of the electrodes and continues in the tissue region between the electrodes. As has been shown by histological analysis, the stimulated growth of the bone follows the gradients of the electric potential. If the electrodes extend through the gap of a fracture, generally perpendicular to the fracture surfaces, \"bone bridges\" form under the influence of the induced potentials and they can quickly take over the supporting fuction of the osteosynthetic elements. As soon as the load carrying capacity of the newly formed bone is sufficient, the whole osteosynthetic part can be removed from the tissue together with the implanted coil and the electrodes.\nFor bone fracture treatment and bone regeneration experience has shown that a few discrete electrodes, for example two to four electrodes can be used which owing to locally concentrated electric potentials stimulate growth which starts and accelerates regeneration of the tissue. The purpose of the electrodes is a formation of a load carrying bone bridge where the bone is defective. The further conduct of regeneration and consolidation of the bone then takes place under the influence of functional loading in accordance with Wolff's Law.\nAs regards an implant which is to remain in the tissue for extended periods of time, if possible for the whole of the life of the patient, the requirements applying are in many ways different and they are in effect very much more exacting; the implant is to take over the functions of missing bone components, the support and bearing of an artificial joint or the anchoring of a tooth in the alveolus of the jaw bone.\nMetallic implants have not been found suitable for this purpose, since the metal is attacked and not accepted by the tissue. Furthermore ceramic materials developed in more recent times for implantation purposes on the basis of aluminum oxide, quartz, silicon nitride and ceramic compound materials still leave much to be desired. While they have a sufficient mechanical strength and there is a freedom from any substantial attack by body fluids ceramic implant materials have not successfully overcome the problem, as was the case of previously employed metallic implants, of the still insufficient degree of biological tolerance. Even in the case of experiments on animals the process of embedding of, for example, a ceramic joint prosthesis lasts approximately three months despite the high regeneration capabilities of animals. During this time the joint must not be used, that is to say the extremity must be kept still. In the case of accident casualities or aged persons the implant does not become embedded or only does so in a substantially longer period of time. However even periods of one quarter of a year in which the patient must keep the fracture still lead to substantial additional health problems.\nAttempts have been made to shorten these times in which the patient must keep still by a special shaping of the implants, for example by the use of a threaded or support rib design, which anchors the implant in the bone with or without the use of a cement. While implants of this type are capable of taking up loads sooner, a functional contact of the implant with the vital tissue is restricted in the case of ceramic prosthesis designed upon these lines to the times applying for metal prosthesis, that is to say on average to approximately 6 to 7 years. An insiduous progressive encapsulation by connective tissue of the implant which is experienced as a foreign body thus leads in the course of time to a painful loosening of the prosthesis. It then becomes necessary to remove the prosthesis and carry out a new implantation, whose life expectancy will, however, be substantially reduced. With the speedy increase in degenerative and accident damage to the human skeleton, which has to be repaired by the implantation of bone prosthesis (at a rate of presently approximately 1500 cases daily) a medicinal and economic problem has been developing on a substantial scale."}
{"text": "Such ejectors are known for instance from German Patent DE 43 02 951 C1. Such ejectors operate on the Venturi principle. The filtered, lubricant-free compressed air flows into the ejector via a connection neck and reaches the propellant nozzle, where the flow velocity of the air acting as a propellant gas is exceeded to supersonic speed in the cross-sectional reduction. After emerging from the propellant nozzle, the air expands and flows via a receiver nozzle into the expansion portion and from there out into the open, optionally via a muffler. In the process, a negative pressure is created in the chamber surrounding the propellant nozzle, and this causes air to be aspirated via the suction connection. The aspirated air and the propellant gas introduced into the ejector emerge jointly into the open via the expansion portion.\nEjectors of this kind have the advantage over vacuum pumps of not having any rotating parts and therefore of requiring little maintenance. They also have low wear. They are also explosion-proof, since as a rule they are pneumatic in nature. Furthermore, they are simple in design and can be installed in an arbitrary position.\nAlso, they develop no heat and can be turned on and off at any time, which saves energy. In addition, the vacuum can be built up faster because of short line lengths between suction cups and the ejector. Finally, the compact design, low weight, and the possibility of combining multiple functions in one device, along with economies in the area of construction, work preparation, procurement, mechanical machining, installation, startup, and spare parts inventory, all play important roles.\nFor fast decay of the vacuum and ejection of the workpiece, in special ejectors the suction connection can be occupied by compressed air, so that an overpressure is applied to the workpiece instead of a negative pressure. However, it has been found that above all when there are high resistances in the suction connection, especially with long connecting lines between the ejector and a suction cup, for instance, engaging the workpiece, the compressed air does not escape, or only a negligible portion of it escapes, via the suction connection and instead it escapes through the receiver nozzle via the suction conduit. In other words, a pneumatic short circuit is created. The attempt has been made to overcome this disadvantage with check valves, but this has led to unsatisfactory results."}
{"text": "The present invention relates to the inspection of components and, more particularly, to an improved method for automatically inspecting gas turbine engine components having complex geometric shapes using eddy current techniques.\nEddy current inspection is a commonly used technique for detecting discontinuities or flaws in the surface of a gas turbine engine component. Eddy current techniques are based on the principle of electromagnetic induction in which eddy currents are induced within the material under inspection. Eddy currents are induced in a test specimen by alternating magnetic fields created in the coil of an eddy current probe when the probe is moved into proximity with the component under test. Changes in the flow of eddy currents are caused by the presence of a discontinuity or a crack in the test specimen. The altered eddy currents produce a secondary field which is received by the eddy current probe coil or by a sensor coil in the eddy current probe which converts the altered secondary magnetic field to an electrical signal which may be recorded on a strip chart. An eddy current machine operator may then detect and size flaws by monitoring and reading the signals recorded on the strip chart. Flaws or defects are detected if the electrical signal exceeds a predetermined voltage threshold.\nPresent eddy current inspection methods work satisfactorily when the components under inspection have simple geometrical shapes, such as holes, flat plates or the like. However, when the component under test has a complex geometrical shape, such as the dovetail slots of a high pressure or low pressure turbine disk, fan disk, high pressure compressor disk, teeth of a gear or the like, the complex geometry of these components such as edges, transitions between convex, concave and flat regions, produces contributions to the eddy current signals which make it difficult to distinguish between defects and non-defects.\nA presently used method of detecting cracks or defects in a complex gas turbine engine component involves scanning a portion of the surface of the complex component with an eddy current probe and converting the received eddy current signals to a two-dimensional digital image. The two-dimensional image is then matched or compared to a chosen template by known image analysis techniques, such as convolution, to detect a defect or flaw. The template is chosen according to the active region of the eddy current probe as well as the size and shape of the defects desired to be detected. The image analysis matching technique will detect a defect only if the size and shape of the defect correspond substantially to that represented by the chosen template. Different templates must, therefore, be used for detecting defects of different sizes and shapes. In order for this detection process to be independent of defect size and shape, too numerous a number of templates would have to be compared to the two-dimensional image to detect all possible flaw sizes and shapes that may be present; such a process would be computationally impractical. Thus, the template matching technique can result in erroneous results if the template does not correspond substantially to the defect and the technique is inefficient because of the large catalog of templates which must be compared to the two-dimensional image."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit device in which the malfunction of a motor drive circuit due to a parasitic effect is prevented.\n2. Description of the Related Art\nConventional three-phase motor drivers have three parallel-connected transistors (Tr1-Tr2, Tr3-Tr4, and Tr5-Tr6) each connected in series between direct-current power supplies VCC and GND. Output terminals taken from between Tr1-Tr2, Tr3-Tr4, and Tr5-Tr6 are connected to a motor M. Thus, normal/back electromotive force is generated as the motor rotates/stops. It is proposed in the art to provide a structure in which a protective diode is connected between the collector and the emitter of the transistors to dissipate the electromotive force to a fixed potential, thereby protecting the interior of an IC including the series transistors (for example, refer to Patent Document: JP-A-6-104459, pp. 13-14, FIGS. 16-17).\nAlso a conventional normal/reverse rotation control circuit of a DC motor is disclosed (for example, refer to Nonpatent Document: Hirobumi Miura; “Mechatronics,” Ohmsha, Ltd., pp. 204-205\nIn the conventional semiconductor integrated circuit devices, for example, during ON to OFF transition of a driver device, backward electromotive force (hereinafter, referred to as back electromotive force) is generated from a motor. Thus a negative potential is applied to the collector region of the driver device for driving the motor by the back electromotive force. Accordingly, free carriers (electrons) are generated from a PN junction between the emitter region and the base region of a parasitic transistor constructed of a driver device, a substrate, and a control device. The free carriers (electrons) flow from an island region having the driver device to the other island regions via the substrate. Particularly, when the free carriers (electrons) flow into the control device for controlling the driver device, the control device fails to function normally. This causes the driver device which should be in OFF action to make ON action to send an incorrect signal, thus interfering with the normal operation of the motor."}
{"text": "The present invention relates to a microorganism capable of producing L-lysine and a method for producing L-lysine thereby. More specifically, the present invention relates to a strain possessing resistance to arginine analogues and other analogues besides the characteristics known as necessary for the production of L-lysine, and a method for producing L-lysine comprising culturing the strain in a nutrient culture medium and recovering L-lysine from the culture broth.\nL-Lysine, which is one of essential amino acids, has been used as an animal feed supplement because of its ability to improve the quality of feed by increasing the absorption of other amino acids, and used as a food supplement as well. L-Lysine has been used also in medicine, particularly as ingredients of infusion solutions.\nL-Lysine can be produced by fermentation using an auxotrophic mutant and/or regulatory mutant, or by enzymatic conversion which comprises treating DL-.alpha.-amino-.beta.-caprolactam (hereinafter referred to as ACL) with L-ACL-hydrolase and ACL-racemase. For commercial production of L-lysine, however, fermentation processes have been more predominantly employed. For example, Japanese Patent Laid-Open No. 82-9797, and Japanese Patent Publication Nos. 82-14839 and 86-35840 disclose the methods using strains having characteristics of vitamin B.sub.1, pantothenic acid and/or biotin requirement for growth; resistance to lysine, threonine and/or isoleucine analogues; resistance to one or more antibiotics e.g. bacitracin, penicillin G and/or polymyxin."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for assembling a camera module for a vehicle.\n2. Description of Related Art\nCamera modules are being increasingly used these days in vehicles, for example for use in night vision systems or as lane departure warning (LDW) systems. Particular demands in terms of robustness are made on such cameras with regard to their installation in the vehicle, given the vibrations that unavoidably occur during operation.\nTwo fundamentally different approaches are known for the assembly of camera modules for vehicles: An imager module is used which encompasses an image sensor chip, a basic housing, and a lens, the lens being produced and tested as an independent module. Imaging modules of this kind can be attached to the circuit board, for example by thermal joining, e.g., using soldering methods, or by way of a cable end having a plug. The circuit board, completed in this fashion, is then installed into the housing, which must have a corresponding receptacle for the lens. A disadvantage of this procedure is the fact that imager modules having chip and lens dimensions suitable for motor vehicle applications cannot be reliably produced using common soldering methods (e.g., reflow methods). Because of the high heat capacity, poor solder connection points can form, or the optics can be damaged. On the other hand, the plug contacts known as alternative solutions are not fault-free over the long term in vehicles. A further problem that has emerged here is sealing of the entire camera module in the context of a circuit board plus imager module configuration, since mechanical stresses acting on the imager module often occur; this can result in subsequent damage.\nA further assembly concept is based on configuration of the circuit board with an independently packaged image sensor; this typically involves a ceramic housing having a glass cover. The sensitive pixel array can thereby be sufficiently protected from particles. Such housings can moreover be processed using reflow processes. The lens can be focused either in a lens holder mounted on the circuit board, or in the housing shell of the camera module. This approach, however, exhibits a very long tolerance chain because of the larger number of individual components requiring installation; after assembly this can result, in the least favorable case, in inaccurately operating camera modules due to addition of the greatest tolerances.\nPublished German patent document DE 10 2004 001 698 discloses an optical module. According to this approach, an optical module encompasses an image acquisition device made up of optical and electronic components, the components of the image acquisition device being disposed together on a carrier plate. The components of the image acquisition device are disposed on a first principal surface of the carrier plate, an optical system associated with the image acquisition device being disposed on a second principal surface of the carrier plate. The carrier plate and the optical system are embedded in a sealing compound. The optical module encompasses a holding element for securing the optical module to a vehicle part or the like. The holding element is sealed to the optical module, and the optical module encompasses a shield against electromagnetic interference radiation. According to this approach, the shield can also be embedded in the sealing compound. The shield is preferably embodied as a net or lattice embedded in the sealing compound, or can be made up of particles embedded in the sealing compound."}
{"text": "The invention relates to a method and a plant for winding a ship into a coil and applies in particular to the winding of a steel band after hot rolling.\nWhen exiting out of a hot band train, the steel band is carried generally on a roller table, towards one or several coiling machines situated at a certain distance downstream of the roller table, generally under it.\nUpstream of the coiling machines, the band is driven by a clamping device, followed by guiding means which direct the band towards a rotary chuck whereon the band, driven by friction, is wound in superimposed turns.\nFor winding a new band, the head of the band is applied on the external face of the rotary chuck in order to be driven by friction downstream. Then, the band is guided by a curved bending plate so as to be bent while remaining applied on the chuck, along an angular sector which increases gradually, which enhances the friction driving effect. After one revolution, the head of the band passes below the internal face of the hand to form a second turn and so on. In order to keep the tension necessary for winding, the band is held upstream by a device including, usually, a deflecting roll and a pinch roll driven into rotation at an angular velocity slightly smaller than that of the chuck, wherein the band may thus be stretched along a theoretical winding plane tangent, upstream, to the deflecting roll, along a generatrix thereof, and, downstream, to the coil being winding. Usually, the chuck is retractable so as to reduce in diameter to enable the removal of the coil, once it has reached the requested diameter.\nThe document U.S. Pat. No. 2,918,226, for instance, describes a winding plant of such type including, downstream of a deflecting roll placed at the outlet of a roller table, a rotary chuck arranged inside a pit and associated with two winding assemblies, each comprising an application roll of the band on the cylindrical external face of the chuck and a curved bending plate for guiding the band during the winding process. In the arrangement of document U.S. Pat. No. 2,918,226, the assembly is placed on a supporting chassis mounted so that it is capable of being pivoted around an axis parallel to the axis of the chuck and actuated by a jack, so as to move between a spaced apart position and a working position for which the bending plate is substantially parallel to the external face of the chuck. On the other hand, the application roll is mounted itself on an auxiliary chassis hinged on the supporting chassis of the bending plate and pushed towards the chuck by a spring so as to exert an elastic force to apply the band on the chuck, along a pushing direction running through a back-up generatrix of the external face of the chuck and the axis thereof.\nMoreover, the application roll can be driven into rotation around its axis in a reverse direction relative to the direction of rotation of the chuck, in order to apply, to the external face of the band, an additional driving friction force which is added to the main driving friction force exerted by the chuck on the internal face of the band.\nAt the beginning of the winding process, the head of the band may rest on the bending plate situated downstream of the application roll and is guided in the annular gap between the bending plate and the external face of the chuck so as to be winded around the chuck. After one revolution, the head of the band may run below the hand so as to form a second turn, the application roll, pushed by a spring, being capable to move apart slightly.\nIn a more developed embodiment described in document U.S. Pat. No. 3,587,274, an additional roll is mounted, upstream of the application roll and of the chuck, on a chassis hinged around an axis, in order to exert a pre-bending torque on the band.\nUntil now, such devices were satisfactory since they were used for winding steel bands with a relative small thickness, for instance ranging between 2 and 12 mm.\nFor such thicknesses, indeed, the band bends easily and may be guided in the annular gap provided between the chuck and the bending plate, so as to be applied on the cylindrical face of the chuck in order to be driven by friction.\nHowever, it is now contemplated to use devices of such type for winding hot bands with increased thicknesses, up to for example 25 mm. Moreover, one may have to wind products with high elastic limits, for instance of the order of 370 MPa at the temperature of use.\nStill, it has appeared that, for relatively high thicknesses, the usual arrangement of the type represented on FIG. 1 does not enable sufficient bending of the band, in particular at the beginning of its winding process.\nAt that time, indeed, the band is applied only on a small angular sector of the chuck, and the main driving force exerted by the chuck on the internal face of the band depends essentially of the pressure exerted by the application roll of the band to the chuck.\nHowever, in the previous arrangements disclosed by the documents mentioned above, this pushing force, exerted by a spring system or an elastic stop, would be insufficient to enable the winding of a hot band which is relatively thick and rigid."}
{"text": "Glypican-3 is an oncofetal antigen that belongs to the glypican family of glycosyl-phosphatidylinositol-anchored heparin sulfate proteoglycans. Glypicans regulate the activity of several growth factors including Wnts, Hedgehogs, bone morphogenic proteins and fibroblast growth factors (FGFs) (Filmus et al. FEBS J. 2013, 280:2471-2476). Glypicans are characterized by a covalent linkage to complex polysaccharide chains called heparinsulphate glycosaminoglycans. Glypicans are involved in cell signaling at the cellular-extracellular matrix interface. (Sasisekharan et al., Nature ReviewslCancer, Volume 2 (2002). To date, six distinct members of the human glypican family have been identified. Cell membrane-bound glypican-3 is composed of two subunits, linked by one or more disulfide bonds.\nGlypican-3 is expressed in fetal liver and placenta during development and is down-regulated or silenced in normal adult tissues. Mutations and depletions in the glypican-3 gene are responsible for the Simpson-Golabi-Behmel or Simpson dysmorphia syndrome in humans. Glypican-3 is expressed in various cancers and, in particular, hepatocellular carcinoma (“HCC”), melanoma, Wilm's tumor, and hepatoblastoma. (Jakubovic and Jothy; Ex. Mol. Path. 82:184-189 (2007); Nakatsura and Nishimura, Biodrugs 19(2):71-77 (2005). The cell surface form of Glypican-3 is highly expressed in HCC (>50%) and in other cancers including squamous lung cancer (approximately 25%).\nGlypican-3 promotes tumor growth in vitro and in vivo by stimulating canonical Wnt signaling which induces the cytosolic accumulation and nuclear translocation of the transcription factor β-catenin (Filmus, supra). It has been shown that GPC3 can form a complex with several Wnts (Capurro et al., Cancer Res., 2005, 65:6245-6254), and Frizzleds, the signaling receptor for Wnts (Filmus, et al., Genome Biol., 2008, 9:224).\nHCC is the third leading cause of cancer-related deaths worldwide. Each year, HCC accounts for about 1 million deaths. (Nakatsura and Nishimura, Biodrugs 19(2):71-77 (2005).) Hepatitis B virus, hepatitis C virus, and chronic heavy alcohol use leading to cirrhosis of the liver remain the most common causes of HCC. Its incidence has increased dramatically in the United States because of the spread of hepatitis C virus infection and is expected to increase for the next 2 decades. HCC is treated primarily by liver transplantation or tumor resection. Patient prognosis is dependent on both the underlying liver function and the stage at which the tumor is diagnosed. (Parikh and Hyman, Am J. Med. 120(3):194-202 (2007).) Effective HCC treatment strategies are needed."}
{"text": "The invention relates to a new and improved method of fabricating hollow cylinders of machines which are subjected to a high working pressure and pronounced erosion effects. Such machine parts are subjected to extremely great wear by virtue of resulting internal pressures in the order of magnitude of 1,000 to 2,000 bar, the frequently high corrosion loads, and last but not least by virtue of the fillers used in processing plastics. It has been found that certain fillers such as, for instance, glass fibers can reduce the service life of such parts to as much as 1/10 of the original value. Since in injection molding machines approximately 0.25 mm constitutes the maximum permissible clearance between worm and cylinder, the significance of the wear factor can be appreciated and, accordingly, on the part of the machine manufacturer considerable efforts have already been made to solve this technological problem.\nThe solution presently most employed is the use of nitrided steels, which render possible the application of a wear-inhibiting nitrided layer to the inner cylinder surface by means of the methods used therefor (bath nitriding, gas nitriding, ion nitriding). However, there results the drawback that during the treatment a pronounced distortion or, respectively, a change in dimensions frequently occurs and the thin nitrided layer is non-uniformly taken off already during the mechanical finishing of the cylinder. Since the cylinder material has relatively low strength, the wear following the machining of such weak locations progresses extremely rapidly. For the same reason the use of case hardened layers, also, has not been sufficiently satisfactory. Attempts have also already been made to provide hot-work steels with hard surface layers, for instance by treating with boron, flame spraying or by means of the so-called CVD-methods (chemical vapor-phase deposition). Therein, however, problems have occurred with respect to dimensional changes due to the required high coating temperatures and to chipping-off of the hard layers during operation.\nAn interesting although expensive solution is represented by the manufacture of so-called \"bimetallic cylinders\". These are understood to be primarily used for machines having larger worm diameters. Specially developed coating alloys are applied to a pre-fabricated part thereof to form an approximately 1.5 to 4 mm thick layer thereon in a hermetically sealed furnace at high rotational speeds needed until cooling (see, for example \"Industrieanzeiger\" Nr. 60, dated July 20, 1973, pages 1403/4). The respective centrifugal casting, however, it not easy to control; also, the brittleness of the bimetallic layers can lead to problems."}
{"text": "It is conventional to provide a substantially horizontal, underground passageway for the accommodation of a conduit having a cylindrical wall provided with openings through which fluids such as water or other liquids may pass to enable earth above the passageway to be dewatered or drained. Examples of apparatus for accomplishing this objective are disclosed in Wagner et al U.S. Pat. No. 4,543,013, issued Sep. 24, 1985 and Cherrington et al U.S. Pat. No. 4,679,637 issued Jul. 14, 1987.\nIt also is conventional to install in an underground passageway, after the formation thereof, a conduit composed of a plurality of pipe sections at least some of which have a wall provided with openings through which water or other fluids may pass. Examples of such conduit-forming pipe sections are disclosed in Lange U.S. Pat. No. 6,073,659 issued Jun. 13, 2000 and Bloomfield et al U.S. Pat. No. 6,561,732 issued May 13, 2003.\nThe disclosures of the above identified patents are incorporated herein by reference.\nThe underground passageway conventionally is formed by a drill member secured at one end of a rotatable drill stem composed of a string of tubular rod sections joined to one another in such manner as to enable all of such sections and the drill member to be rotated simultaneously and moved linearly in a direction to form the underground passageway. Following forming of the passageway the drill member conventionally will occupy an accessible position enabling it to be removed from the drill stem and the drill stem then to be coupled to a reamer which can be drawn rearwardly through the passageway so as to enlarge the passageway. The reamer may be coupled to one end of a string of conduit-forming pipe sections which will be drawn through the reamed passageway in response to rearward movement of the drill stem and reamer. When the pipe string reaches the point at which the forming of the passageway commenced, it may be detached from the drill stem and coupled to a pump so as to enable whatever liquid is in the underground conduit to be withdrawn therefrom. In this manner a swampy area above the conduit may be drained. Further, contamination which may result from the discharge of toxic materials to land areas may be drained from such areas. In addition, fluids containing bioremedial substances may be pumped into the conduit and discharged through the openings thereof into the area encircling the conduit so as to detoxify such area.\nThere are several objectionable characteristics of the conventional apparatus referred to above. Among them are the necessity of making at least two passes underground, one of which is in the initial formation of the passageway and the other of which is the reaming of the passageway to an appropriate size to accommodate conduit and installing the conduit. Another objectionable characteristic is the difficulty of disconnecting the drill stem from the drill member or other device rotated by the drill stem, including in some instances the necessity of extending the drill member to the surface or excavating a pit to obtain access to the drill member in its underground position. A further objectionable characteristic of conventional apparatus is that the openings in the conduit may not be of such size as to accommodate fluids of different viscosities with optimum flow rates.\nA principal object of the invention is to provide methods and apparatus which overcome the objections referred to above."}
{"text": "1. Field of the Invention\nThe present invention relates to an image processing device that performs white balance control, a method of controlling the same, and an image pickup apparatus, and more particularly to white balance processing performed by the image processing device.\n2. Description of the Related Art\nIn recent years, in automatic white balance processing performed by an image processing device, such as an image pickup apparatus, a so-called TTL (Through The Lens) method is widely used. In automatic white balance processing using the TTL method, the color of light from a light source is estimated according to an image obtained by photographing. Further, to calculate a white balance correction value (hereinafter referred to as a WB correction value), it is necessary to accurately distinguish, on the image, between the color of light from the light source and the color of the object.\nFor example, there has been proposed an image pickup apparatus that divides an image into a plurality of areas (block areas), and applies color information obtained from each divided area to detection frames each indicative of a range of color distribution corresponding to the light source type (see Japanese Patent Laid-Open Publication No. 2003-163944). In this image pickup apparatus, the light source is determined, on a detection frame-by-detection frame basis, according to the number of applicable areas, the luminance level of an object, and an infrared light amount, and white balance control is performed based on the result of the determination.\nFurther, there has been proposed an image pickup apparatus that is configured to perform proper white balance control with respect to a green object included in a white extraction range, on which an action for converting green to an achromatic color works (e.g. in a case where green trees or grasses are photographed in a dark place, such as a place in the shade) (see Japanese Patent Laid-Open Publication No. 2006-174281). In this image pickup apparatus, white balance is controlled by determining, based on distribution of data on a three-dimensional coordinate system having an illuminance, an amount of green, and an R gain, as coordinate axes, whether an image has been photographed indoors using indoor fluorescent lamps as a light source, or has been photographed outdoors in the sun/shade.\nHowever, in the image pickup apparatuses described in Japanese Patent Laid-Open Publication No. 2003-163944 and Japanese Patent Laid-Open Publication No. 2006-174281, even under the light from the same light source, if the color of an object in an image changes, it is difficult to perform white balance control using the same WB correction value.\nFor example, the image pickup apparatus described in Japanese Patent Laid-Open Publication No. 2003-163944 considers using an intermediate value of the WB correction values associated with the respective light sources, by changing a gain value according to an evaluation value of each light source. However, in Japanese Patent Laid-Open Publication No. 2003-163944, since the light source is detected based on the color distribution, if the color of an object in the image changes, it is difficult to stably continue application of the WB correction value.\nFurther, the image pickup apparatus described in Japanese Patent Laid-Open Publication No. 2006-174281 considers stepwise changing the white balance control according to a distance from a light source identification boundary on a three-dimensional coordinate plane. However, also in Japanese Patent Laid-Open Publication No. 2006-174281, if the color of an object changes, it is difficult to stably continue application of the WB correction value, similarly to the image pickup apparatus described in Japanese Patent Laid-Open Publication No. 2003-163944."}
{"text": "This section introduces aspects that may help facilitate a better understanding of the disclosure. Accordingly, these statements are to be read in this light and are not to be understood as admissions about what is or is not prior art.\nChronic pain is a major health concern that costs the US more than $635 billion per year (Gaskin and Richard, (2012) J. Pain 13:715-724). In addition to the financial impact, patients with chronic pain suffer extreme physical, emotional, and social burdens. For example, individuals often become socially isolated and confined to home as a result of their chronic pain that is not well-controlled by today's available treatments. The drugs used for the management of chronic pain include opioid analgesics, neuronal stabilizers such as anticonvulsants, and antidepressants. Opioids are the most widely used, and a recent NIH report indicates that there are significant problems associated with long-term opioid therapy for chronic pain (Volkow and McLellan (2016) N Engl J Med 374:1253-1263.). None of the agents provide sufficient relief to allow patients to return to their normal activity level. Moreover, current pharmaceutical industry has retreated from studying novel pain therapeutics due to the enormous risk (Skolnick and Volkow (2016) Neuron 92:294-297.). These observations indicate an essential need to identify new agents acting on unique targets in the war on chronic pain.\nAdenylyl cyclases are enzymes that catalyze the production of cyclic adenosine monophosphate, an important biological messenger. There are nine different membrane-bound human adenylyl cyclase isoforms, each with slightly different attributes. In particular, adenylyl cyclase 1 (AC1) is highly expressed in the hippocampus as well as in regions of the brain associated with pain. AC1 knock out animals show reduced neuropathic and inflammatory pain. Additionally, AC1 knock out mice show less reward when given opioids and show reduced symptoms of opioid dependence during withdrawal. Additional reports suggest that AC1 inhibition may also provide a useful therapeutic intervention for alcohol use disorder and autism (Bosse K E et al., J. Pharmacol. Exp. Ther. 2017, 363 (2) 148-155; Sethna F., et al. Nat. Commun. 2017, 8, 14359).\nUnfortunately, until now, the selective inhibition of ACs has not been achieved, and simultaneous inhibition of multiple adenylyl cyclase isoforms would likely result in significant adverse effects. There are unmet needs for better and safer medications targeting adenylyl cyclases for various therapeutic uses."}
{"text": "1. Field of the Invention\nThe invention relates to methods and systems for detecting access points, and particularly to methods and systems for detecting rogue access points and a device for identifying rogue access points.\n2. Related Art\nWhen an organization such as a company uses a traditional wired network communication and computer system, it is generally not easy for network intruders to gain access into the company's network. For instance, the network intruders may need to gain physical access inside the company's building, and then find an Ethernet port from where they can electronically obtain confidential information of the company. This physical boundary provides a certain level of protection against network intruders.\nWhen a company uses a wireless network, the network intruders do not necessarily have to physically get inside the company's building in order to access the company's network. The network intruders only need to be present within the scope of the wireless network, and then gain access to the company's confidential information through a wireless connection.\nOne form of wireless security breach occurs when network intruders put an access point (AP) near to or in a company's building. An employee of the company may unknowingly connect the unauthorized access point to the company's network. The network intruders can then try to gain access to the employee's confidential information via the unauthorized access point. This kind of unauthorized access point is commonly known as a “rogue access point” or “rogue AP.” Therefore, a method for detecting rogue access points is needed in order to protect an organization's confidential information."}
{"text": "This invention relates to a method for preparing a therapeutic agent that, upon ingestion decreases intestinal cholesterol absorption in man and, specifically inhibits or decreases intestinal cholesterol absorption by inhibiting the pancreatic cholesterol esterase catalyzed hydrolysis of naturally occurring and ingested cholesterol esters and by inhibiting the cholesterol esterase facilitated uptake of free cholesterol.\nThe invention is based upon the discovery that pancreatic cholesterol esterase is an important contributor to overall dietary cholesterol absorption because (1) cholesterol derived from cholesterol esters is preferentially absorbed compared to free-cholesterol; (2) cholesterol esterase enhances the absorption of free cholesterol and (3) dietary cholesterol and/or cholesterol esters induce the mRNA and enzymatic activity of cholesterol esterase in the pancreas in a newly discovered intestinal-pancreatic cycle for the absorption of cholesterol. U.S. Pat. Nos. 5,173,408 and 5,063,210 describe the importance of cholesterol esterase in the dietary uptake of cholesterol and also disclose methods for inhibiting cholesterol esterase. Thus, the surprising usefulness of inhibiting cholesterol esterase has demonstrated a new need for potent (Ki less than 5 xcexcM) and safe inhibitors of cholesterol esterase.\nMany physical ailments are attributed at least in part to high levels of serum cholesterol. Atherosclerosis, for example, is a leading cause of death in the United States and high serum cholesterol concentrations are associated with increased risks of fatal artherosclerotic events. The discovery that the cholesterol esterase enzyme plays a role in intestinal cholesterol absorption has led to attempts to attenuate intestinal cholesterol absorption in man by inhibiting the action of the cholesterol esterase enzyme.\nThere is now an important need to develop human pancreatic cholesterol esterase inhibitors, especially those that are not absorbed essentially nondegradable. The pharmacology of various polysaccharides has been investigated. Cook and Cammarata, 1963, Arch. Int.: Pharmacodyn. 144: 1. In particular crude sulfated amylopectin has been taught in U.S. Pat. No. 4,150,110 as an anti-ulcer agent, but its property as a cholesterol esterase inhibitor has not been recognized. Sulfated dextran of low molecular weight has been recognized for use in the treatment of hyperlipemia and as an orally administered anticoagulant, British Patent No. 953,626.\nLow molecular weight sulfated dextran (MDS-T) has been used to reduce serum cholesterol levels in Japan at a dose of 1800 mg/day, Goro et al, 1987, J. Clin. Biochem. Nutr. 2: 55-70 by activating a blood enzyme lipoprotein lipase. The low molecular weight of this bacterial dextran, (7-8 kDa), allows the sulfated dextran to be absorbed by the intestine as demonstrated by carbon-14 labelling studies. Drugs In Japan-(Ethical Drugs, 10th ed. 1986). It was developed for this property of intestinal absorption as indicated by the claim that faster decreases in serum lipids can be obtained by intravenous administration of this agent with clearance of serum lipemia due to activation of plasma lipoprotein lipase. Clearly this route of administration will not lead to effects on inhibiting cholesterol esterase in the intestine. Absorption of MDS can lead to a variety of side effects, most notably, anticoagulant effects that must be monitored. This preparation has not been known to inhibit cholesterol esterase and it is sulfated randomly and at various ring positions.\nMore recently, it has been discovered that crude non-absorbable polysaccharides sulfated at the three position are most effective as inhibitors of cholesterol esterase. U.S. patent application Ser. No. 08/121,369. Useful 3-sulfated polysaccharides may be derived from the synthetic sulfation of polysaccharides from various natural sources including seaweeds.\nMethods for preparing sulfated polysaccharides are also known in the art. For example, U.S. Pat. No. 3,624,069 describes the sulfation of cellulose with a sulfur trioxide/lower n-dialkyl amide sulfation complex. U.S. Pat. No. 4,480,091 describes a process for preparing cellulose sulfate esters in a three step process. Finally, U.S. Pat. No. 4,814,437 describes a method for preparing sulfated polysaccharides by subjecting the polysaccharide to a reducing step prior to sulfation.\nThe present invention is directed to a method for manufacturing high molecular weight 3-sulfated polysaccharides that are essentially non-absorbable and nondegradable in the alimentary tract, and when administered orally, they are useful in decreasing human serum cholesterol and LDL levels by inhibiting human pancreatic cholesterol esterase, now recognized as a key enzyme involved in mediating cholesterol absorption. Thus, following the methods of this invention, greater than 95% of the compound has a molecular weight greater than 75,000 Daltons. The sulfate to monomer ratio is between 1.0 and 3.0 and less that 0.5% (by wt.) of the material will be free sulfate. This very high molecular weight compound can be administered to humans in tablet form, incorporated in a foodstuff, or by any other method that inhibits cholesterol absorption in the alimentary tract."}
{"text": "Reset signals may be used to initialize electrical circuits as a part of various system operations. In some cases, reset signals may be used to force a circuit into a known state for simulation or a defined state desired by a designer or a user.\nA synchronous reset signal is often used to affect or reset state of flip-flops on the active edge of an associated clock signal. Thus, the synchronous reset signal should have a duration that is sufficient to be captured at the active clock edge. While a synchronous reset signal provides some immunity to glitches, unless they occur right at the active edge of an associated clock signal, it may not be very useful in power saving schemes where the clock signal is gated.\nAn asynchronous reset signal may be used to affect the state of flip-flops and other circuitry asynchronously, without regard to the state of the clock signal. This may be useful in high speed circuits, as the data path becomes independent of reset signal. However, an asynchronous reset signal may not offer the same level of immunity to glitches that is provided by the synchronous reset signal, contributing to spurious circuit reset operations."}
{"text": "1. Field of the Invention\nThe present invention relates to an automated system for handling a load on a pallet and more particularly to such a system which transfers a load from a first pallet to a second pallet, with or without a slip sheet, or to a slip sheet only.\n2. Description of the Related Art\nManufacturers of products, such as powdered laundry soap, typically stack a preselected number of packages containing the product on a wood pallet for shipping. A load of product supported by a pallet is referred to as a palletized load. Such stacked products are often wrapped with a stretchable polymeric film for maintaining the stacked packages as a unit on the shipping pallet.\nSometimes the manufacturer initially stacks the packages on a captive pallet, so called because it is a high quality, durable pallet which is not shipped with a palletized load, but rather is retained by the manufacturer for repeated use. The captive pallet is loaded with product immediately after the product is packaged. The captive pallet is used for transporting the load from the manufacturing area to a warehouse or the like and from there to the point of shipment.\nImmediately prior to shipping, the palletized load is transferred from the captive pallet to a cheaper, less durable shipping pallet. The captive pallet is then re-used by the manufacturer for transferring product from place to place as described above and the shipping pallet containing the product is loaded on a carrier for transporting the goods.\nU.S. Pat. No. 3,123,232 to Postlewaite discloses a pallet load transfer mechanism in which a first pallet and a load of product stacked thereon are sandwiched between a pair of opposed clamps and pivoted over 90.degree. to permit removal of the first pallet and insertion of a second pallet beneath the load. The Postlewaite mechanism is complex and requires high-power motors for driving the mechanism.\nU.S. Pat. No. 3,645,409 to Sinclair discloses a load transfer and pallet stacker which is used to transfer a load on a pallet to the forks of a forklift. Although a load so transferred may thereafter be deposited upon a second pallet by the forklift, it would be desirable to fully automate the process of transferring a load from a captive pallet to a shipping pallet.\nThe present invention comprises a substantially horizontal frame. A substantially planar push plate is spaced above and parallel to the frame and is mounted thereon for horizontal movement along a predetermined axis. A loaded captive pallet staged onto the frame and into the travel path of the plate.\nIn one aspect of the invention, a retainer restrains the load while the plate is moved against the captive pallet thereby pushing the pallet from beneath the load and transferring the load at the plate. A shipping pallet is staged onto the frame beneath the push plate. Continued plate movement withdraws the plate from beneath the load and transfers it from the plate to the shipping pallet. A slip sheet may be provided on the shipping pallet if desired.\nIn another aspect of the invention, a slip sheet is interposed between the load and the captive pallet. The captive pallet is staged onto the frame and onto the travel path of the plate. A slip sheet gripping assembly grips one side of the slip sheet while the plate moves against the pallet and pushes the same from beneath the slip sheet thereby transferring the slip sheet and the load thereon to the plate. A shipping pallet is staged onto the frame beneath the push plate so that continued plate movement withdraws the plate from beneath the slip sheet thereby transferring the slip sheet and the load supported thereon from the plate to the shipping pallet.\nIn still another aspect of the invention, the slip sheet gripping assembly operates as described above to remove the slip sheet and the load supported thereon from the captive pallet without transferring the same to another pallet.\nIt is a general object of the present invention to provide a pallet load transfer method and apparatus which overcomes the disadvantages associated with such prior art methods and apparatus.\nIt is another object of the present invention to provide such a method and apparatus which operates to transfer a load with or without a slip sheet on a first pallet to a second pallet.\nIt is a more specific object of the present invention to provide such a method and apparatus which operates to transfer a load without a slip sheet on a first pallet to a second pallet with or without a slip sheet.\nIt is still another object of the present invention to provide such a method and apparatus to transfer a load with or without a slip sheet on a first pallet to a slip sheet only.\nIt is yet another object of the present invention to provide such a method and apparatus which is fully automated and which performs such transfers at a relatively rapid rate.\nIt is still another object of the present invention to provide such a method and apparatus to transfer a load with a slip sheet on a first pallet to a second pallet with a slip sheet or to a slip sheet only without touching the sides of the load.\nThe foregoing and other objects, features and advantages of the invention will become more readily apparent from the following detailed description of a preferred embodiment which proceeds with reference to the drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a color linear image sensor apparatus, and more specifically to a color linear image sensor apparatus comprising three arrays of photosensor cells.\n2. Description of Related Art\nReferring to FIG. 1, there is shown one conventional color linear image sensor apparatus comprising three arrays of photosensor cells. This conventional color linear image sensor apparatus includes three arrays 1R, 1G and 1B of photoelectric conversion cells, namely, photocells, and a CCD register 2R-1, 2R-2, . . . located at each side of each of these photocell arrays 1R, 1G and 1B. In this example, these photocell arrays 1R, 1G and 1B have a peak sensitivity in red, in green and in blue, respectively, so that the apparatus functions as a color linear image sensor. With this arrangement, since the three linear image sensors are located in parallel, a distance \"L\" between each pair of adjacent photocell arrays is ordinarily required to be on the order of 100 .mu.m to 200 .mu.m.\nBecause of this distance \"L\" between each pair of adjacent photocell arrays, each of the photocell arrays senses a different position on the same subject copy, as will be apparent from FIG. 3 illustrating an optical path diagram in the case of using the color linear image sensor apparatus. Color signal processing needs a color signal for red, a color signal for green and a color signal for blue on the same position. Because of the difference between the positions sensed by the photocell arrays, it was necessary to externally store the color signals of the amount in proportion to the distances \"L\" between the photocell arrays. For example, assuming that the distance \"L\" between each pair of adjacent photocell arrays is 140 .mu.m, a size \"d\" of each photocell is 14 .mu.m, and a signal for each one pixel has information of 10 bits, a memory capacity \"M.sub.1 \" required for each one pixel is expressed as follows: EQU M.sub.1 =10(L+2L)/d=300 (bits)\nFurthermore, assuming that one line is composed of, for example 5,000 pixels, a memory capacity \"M.sub.T \" required for processing the color signals for one line is expressed as follows: EQU M.sub.T =M.sub.1 .times.5000=1.5 (megabits)\nin general, memory for this purpose is required to have a high speed reading operation; therefore, an expensive memory such as a static random access memory (SRAM) has been used. Accordingly, there is a demand for a color linear image sensor apparatus having a shortened distance \"L\" between adjacent photocell arrays, so that a required external memory capacity can be made small.\nIn the above mentioned conventional color linear image sensor apparatus, as will be understood from FIG. 1, there is restriction in shortening the distance between adjacent photocell arrays, because there are wiring conductors for clocks .phi..sub.1 and .phi..sub.2 between the CCD register 2R-2 and the CCD register 2G-1.\nFor example, referring to FIG. 2, there is shown an enlarged layout pattern diagram of a portion \"A\" in FIG. 1.\nIn FIG. 2, the CCD register 2R-2 and the CCD register 2G-1 are located between two photocell arrays 1R and 1G, each of which is formed by arranging a number of photodiodes PD (isolated from each other by a channel stopper 6) in a single line. The CCD register 2R-2 and the CCD register 2G-1 have first transfer gate electrode 11-1 (formed of a first level polysilicon film) and a second transfer gate electrodes 11-2 (formed of a second level polysilicon film), which are common to the CCD register 2R-2 and the CCD register 2G-1. The CCD register 2R-2 and the CCD register 2G-1 are isolated from each other by a channel stopper 6A. Between the photocell array 1R and the CCD register 2R-2, there is formed a transfer gate electrode 12-2 (3R-2) of the second level polysilicon film, and similarly, between the photocell array 1G and the CCD register 2G-1, there is formed a transfer gate electrode 12-1 (3G-1) of the second level polysilicon film. Clock signal wiring conductors 14(.phi..sub.1) and 14(.phi..sub.2), which also function as a light shield, are connected to the first transfer gate electrode 11-1 and the second transfer gate electrode 11-2 through contact holes C1 and C2, respectively.\nThis arrangement was attributable to an idea of forming the clock signal wiring conductors 14(.phi..sub.1) and 14(.phi..sub.2) above the CCD registers. Nevertheless, the channel stopper 6A was required to have a sufficient width to form the contact holes C1 and C2 which connect the clock signal wiring conductors to the first and second gate electrodes, respectively. For example, in the case of designing the layout pattern under the 2 .mu.m rule, assuming that the size of the contact holes is 2 .mu.m.times.3 .mu.m and the alignment margin is 1 .mu.m, the channel stopper 6A is required to have the width of at least 10 .mu.m, and preferably about 12 .mu.m. In this connection, only for the purpose of electrical isolation of the transfer channels, it is sufficient if this width of the channel stopper is at least 2 .mu.m. Therefore, even if the above mentioned design is adopted, the existence of the clock signal wiring conductors is still hindrance in shortening the distance \"L\" between the photocell arrays.\nFurthermore, as one method for shortening the distance \"L\" between the photocell arrays, it might be possible, as shown in FIG. 4, to connect the gate electrodes (11-1A, 11-2A, . . . ) of the CCD register 2 on the channel stopper 6 of the photocell array 1, as adopted in an area image sensor. However, this method requires four-phase, tri-level, pulses .phi..sub.1, .phi..sub.2, .phi..sub.3 and .phi..sub.4 for driving the CCD registers, and therefore, the driving condition becomes complicated."}
{"text": "The invention relates generally to wireless cellular communication systems and more particularly to techniques for automatic handover of an ongoing mobile communication from one cell to another cell within the system, using information transmitted by the base stations or other fixed part system transmitters.\nConventional wireless cellular communication systems include both xe2x80x9coutdoorxe2x80x9d and xe2x80x9cindoorxe2x80x9d systems. A typical outdoor system, such as the AMPS wireless cellular telephone system used in the United States, covers a relatively large geographical area. The system uses radio frequency (RF) signals to establish communication between system base stations and mobile radio units or xe2x80x9cmobiles.xe2x80x9d The mobiles are generally free to move throughout a coverage region of the system, while the base stations are in fixed locations, such as on towers positioned throughout the coverage region. Much of the communication in such an outdoor system takes place when the mobiles are located outside of any particular building or other structure. An indoor wireless system, in contrast, is generally set up within a given building, office complex or other structure or set of structures, and utilizes RF signals to support communications between mobiles and fixed base stations located throughout the interior of the structure(s).\nA common characteristic of both outdoor and indoor wireless communication systems is the use of coverage cells. Each cell is designed to maintain favorable radio characteristics within its boundary. If a mobile engaged in an ongoing call crosses the boundary of a given cell and moves into another cell, the system must be capable of transferring the ongoing call to the other cell. This process is referred to as xe2x80x9chandover,xe2x80x9d and ideally should be implemented so as to cause substantially no disruption in the ongoing call. Wireless systems also use the handover process to overcome or avoid the effects of system disturbances, such as multi-path fading effects. The handover may be from one time or frequency channel to another, or from one cell to another. In many outdoor systems, such as AMPS, the handover process is centrally controlled and thus almost completely autonomous with respect to the mobile. In contrast, the handover process in many indoor systems, including systems based on standards such as DECT in Europe and PWT in the United States, is completely under the control of the mobile and is therefore based primarily on the mobile\"\"s view of the radio spectrum in its own locale.\nThe radio spectrum environment in an indoor system can be particularly harsh, due to random propagation channel effects such as fading and shadowing. The distances between base stations in an indoor system are therefore generally much smaller than those of outdoor systems. From the perspective of the mobile, many base station signals or xe2x80x9cbeaconsxe2x80x9d from other cells can often be detected with excellent signal-to-noise (S/N) ratios at various times and locations as the mobile moves within a given cell. Since the mobile is often in control of any handover actions taken to improve signal quality in an indoor system, the mobile must keep track of all the beacons that it detects, regardless of which cell they originate from, in order to make an intelligent choice from among the available beacons.\nThe mobile in a system with conventional mobile-controlled handover therefore must examine beacons which may or may not be from xe2x80x9cnearest neighborxe2x80x9d or adjacent cells, since it has no way of determining which of the beacons that it detects are from the adjacent cells. Typically, many beacons from both adjacent and non-adjacent cells may be present. For example, it is possible in many indoor environments that a beacon from a non-adjacent cell may momentarily have a sufficient signal strength within the current cell to qualify as a potential handover candidate. An attempted handover to such a candidate will be likely to fail due to changes in multipath or other channel conditions. Even if the attempted handover is successful, changes in channel conditions will tend to necessitate an immediate additional handover. The result in either instance is a slowdown in the overall process of achieving a successful handover, with a corresponding degradation in voice quality. Therefore, despite the fact that many of the beacons that reach a given mobile are not useful in the handover decision, conventional mobile-controlled handover processes require the mobile to consider each beacon in turn in order to weed out the non-useful beacons. This takes time that often results in noticeable degradation of the voice quality in an ongoing call while a handover decision is being made.\nA need therefore exists for improved handover techniques that can speed up the handover process, thereby avoiding the handover delays and corresponding degradation in voice quality associated with conventional techniques.\nThe invention provides methods and apparatus for implementing assisted handover in a wireless cellular communication system. The invention allows a mobile to maintain substantial control of the handover process, while utilizing cell adjacency information stored in the system infrastructure to provide information regarding the most relevant potential handover candidates for that mobile. The cell adjacency information is used to identify the cells adjacent to the cell in which the mobile is currently located, and the potential handover candidates will correspond to beacons generated by these adjacent cells. This allows the mobile to direct its handover attempts to beacons from adjacent cells, while ignoring beacons from non-adjacent cells. The invention thereby substantially reduces handover-related delays and corresponding voice quality degradation in a mobile-controlled handover process.\nIn an illustrative embodiment of the invention, when a mobile conducting an ongoing call in a xe2x80x9ccurrentxe2x80x9d cell determines that a handover to another cell may be needed, the mobile sends a handover assistance request to the base station of the current cell. The base station responds by transmitting to the mobile information regarding a set of potential handover candidates. The potential handover candidates are determined based on stored cell adjacency information which includes identifiers of each of the cells adjacent to the current cell. The adjacency information may be stored in the system base stations after initial system configuration, based on the actual physical layout of the cells, and may be altered during system operation to reflect the continual success or failure of particular attempted handovers.\nThe information regarding the set of potential handover candidates transmitted to the mobile may be in the form of a list of beacons associated with the cells adjacent to the current cell, and may also include corresponding channel identifying information such as frequency and time slot. This transmitted information enables the mobile to focus its handover attempts on handover candidates from adjacent cells, and therefore those candidates which are most likely to result in a successful handover. Alternative embodiments of the invention may transmit the potential handover candidate information from a fixed part transmitter other than a system base station, or may periodically broadcast this information throughout the system rather than transmit it upon receipt of a handover assistance request."}
{"text": "Conventional wearable electronic devices include may include a watch band for attaching the device to a user's wrist. There are varieties of materials and configurations for forming conventional wearable bands. In some cases, it may be desirable to form a band having multiple layers of different materials. However, due to normal or predictable use, a composite or multi-layer band that is constructed using some traditional techniques may be susceptible to failure. The techniques described herein may be used to form a multi-layer or multi-material band for an electronic device that may have a construction that is more robust as compared to some traditional watch bands."}
{"text": "1. Field of the Invention\nThe present invention relates to a plasma reactor and process for the production of hydrogen-rich gas from light hydrocarbons.\n2. Description of the Related Technology\nImproving the efficiency of energy production remains an important technological goal, owing to the significant economic benefits that result in almost every sector of the economy. One potential method for improving the efficiency of energy production is to provide an energy efficient method of converting light hydrocarbons to hydrogen-rich gas, to thereby increase energy production from natural gas.\nPlasma fuel converters such as plasmatrons are known to reform hydrocarbons to produce hydrogen-rich gas. DC arc plasmatrons, for example, are disclosed in U.S. Pat. Nos. 5,425,332 and 5,437,250. DC arc plasmatrons generally operate at low voltage and high current. As a result, these plasmatrons are particularly susceptible to electrode erosion and/or melting. DC arc plasmatrons also require relatively high power inputs of 1 kW or more and relatively high flow rates of coolant to keep the temperature in check.\nOther conventional methods for the conversion of light hydrocarbons to hydrogen-rich gas are generally energy inefficient and, as a result, in many small-scale applications, such as the production of hydrogen for fuel cells, the cost of hydrogen gas made by these methods is not competitive. Thus, there is a need in the art for a more energy efficient process for the conversion of light hydrocarbons to hydrogen-rich gas.\nU.S. Pat. Nos. 5,993,761 and 6,007,742 (Czernichowski et al.) describe processes for the conversion of light hydrocarbons to hydrogen-rich gas using gliding arc electric discharges in the presence of oxygen and, optionally, water. In the process, two electrodes having flat sheet geometry are arranged for arc ignition and subsequent gliding of the arc. The distance between the cathode and anode gradually increases to a point that no longer supports the gliding arc. As a result, the gliding arc disappears at one end of the electrodes, creating pulsed plasma wherein the properties of the plasma change with time. Due to the use of pulsed plasma, the process is relatively unstable over time. Reagents and oxygen are preheated using an external heat source. As a result of the preheating of the reagents and oxygen using an external heat source, the process suffers from poor energy efficiency. A premixed feed gas including hydrocarbons and oxygen is introduced to the reactor located at the central axis of the reactor.\nU.S. Pat. No. 5,887,554 (Cohn et al.) also discloses a system for the production of hydrogen-rich gas from light hydrocarbons. The system includes a plasma fuel converter for receiving hydrocarbon fuel and reforming it into hydrogen-rich gas. The plasma fuel converter can be operated using either pulsed or non-pulsed plasma and can utilize arc or high frequency discharges for plasma generation. Products from the plasma fuel converter are employed to preheat air input to the fuel converter. In one embodiment shown in FIG. 6, residence time in the reactor is increased by providing a centralized anode and a plurality of radial cathodes to thereby cause the arc to glide towards the center of the reactor under the influence of gas flowing in the same direction as the gliding arc.\nU.S. Pat. No. 6,322,757 (Cohn et al.) discloses a plasma fuel converter which employs a centralized electrode and a conductive reactor structure which acts as the second electrode for creation of a plasma discharge. Reagents are fed to the reactor just below the smallest gap between the electrodes and flow in the same direction as the gliding arc to thereby produce hydrogen-rich gas. In alternative embodiments, air and/or fuel are preheated by counter-flow heat exchange with the products of the reforming reaction and fed to the reactor either above or just below the smallest gap between the electrodes.\nAlthough some improvements in the energy efficiency of plasma fuel converters have been achieved, there remains a need for higher energy efficiencies for use of non-equilibrium low temperature plasma."}
{"text": "1. Field of the Invention\nThe present invention relates to a heat sink with fins used for cooling electronic devices, and in particular relates to the heat sink with fins in which heat dissipating fins are jointed on a metal base plate.\n2. Description of the Related Art\nIn order to dissipate the heat generated by a semiconductor chip, which is increasing year by year, there is known a method in which a heat sink is attached in close contact to the semiconductor chip so as to transfer the heat from the semiconductor chip to the heat sink, and dissipate same. A conventional heat sink has a metal flat base plate so as for the semiconductor chip to be closely contacted thereon, and has construction that plural heat dissipating fins of metal thin plates are jointed on the surface opposite to the surface of the base plate on which the semiconductor chip is attached,\nIn the heat sink with fins having the above-mentioned construction, the heat generated by the semiconductor chip transfers from the semiconductor chip to the metal base plate which is closely contacted thereto, and spreads over the base plate, and thus spread heat is dissipated into air or the designated place by the metal heat dissipating fins which are fixed on the surface of the metal base plate. In addition, there is provided a one-piece heat sink made of forged aluminum as one of the conventional heat sinks. However, there is a problem in which the one-piece heat sink made of forged aluminum do not attain enough heat dissipating effect, because required number of heat dissipating fins are not formed due to the technical difficulty by forging to shorten the pitch between the heat dissipating fins for the required level of fin-density.\nIn order to solve the above-mentioned problem and obtain much more excellent heat dissipating effect, there is proposed a method instead of a one-piece heat sink, in which a base plate and heat dissipating fins are separately manufactured and the heat dissipating fins are jointed on a surface of the base plate by means of brazing or mechanical joint.\nThe mechanical joint in which the heat dissipating fins are jointed mechanically on a surface portion of the base plate is superior to the brazing, because manufacturing cost of brazing is higher than that of mechanical joint.\nAs shown in FIG. 4, when the size of the semiconductor chip 30 and the base plate 22 are compared, the semiconductor chip 30 is much smaller than the base plate 22. Therefore it is necessary to spread the heat generated by the small semiconductor chip 30 all over the large base plate 22 by using the base plate made of materials having high heat conductivity such as copper, aluminum or the like. In particular, since the copper is excellent in heat conductivity, the copper is known as a material for the base plate of the heat sink with fins. However, since the copper is heavy and the copper base plate must be in direct contact with the semiconductor chip, there is required some reinforcing construction (for example: applying a jointing device for protecting a circuit board or the like) to reduce the damage that the weight of the copper base plate imposes onto the semiconductor chip.\nFurthermore, electronic devices become remarkably lighter and smaller, which makes it difficult to use the copper base plate because the chance to use such heavy copper base plate is reducing and it becomes difficult to secure the space for installing the copper base plate. On the other hand, when the semiconductor chip becomes more integrated and the processing capacity becomes higher, the heat generated by the semiconductor chip becomes larger. Therefore, further higher heat dissipating effect is required for stable operation of the semiconductor chip or the like.\nAs described above, it becomes difficult to use the copper base plate due to the heavy weight thereof in spite of the excellent heat conductivity. Therefore, when aluminum is to be used for material of the based plate, it is necessary to satisfy the requirement of further increasing heat dissipating effect of the heat sink.\nIn the heat sink having construction that the heat dissipating fins are mechanically fixed on a surface of the base plate, the base plate has the following temperature distribution when the heat dissipation is in the stable condition. More specifically, the temperature is the highest in the central portion of the base plate to which the semiconductor chip is attached in close contact and, the temperature becomes lower as the portion in the base plate moves farther from the semiconductor chip.\nIn addition, the heat dissipating fins which are fixed and thermally connected to a surface of the base plate has the following temperature distribution. More specifically, the temperature is the highest in the heat dissipating fins locating in the central portion of the base plate to which the semiconductor chip is attached in close contact, and the temperature is the lowest in the heat dissipating fins locating in the end portions of the base plate which are far apart from the central portion.\nFurthermore, the heat dissipating fin itself has the following temperature distribution along the height direction. More specifically, the temperature is the highest in the bottom portion of the fin which is fixed and thermally connected to the base plate and, the temperature becomes lower as the portion moves farther from the bottom portion.\nIt proves that the heat is transferred by heat conduction in the heat dissipating fin.\nHowever, it is desirable to minimize difference of temperature across the entire portion of the heat sink with fins, from the heat dissipating efficiency point of view. As described above, lately, the heat generated by the semiconductor chip is increasing, and at the same time, the semiconductor chip is more downsized, which tends to rapidly increase the heat density. Therefore it is urgently required to improve heat dissipating efficiency of the heat sink with fins (for example, increasing density of the heat dissipating fins fixed on the base plate, or using a heat pipe in the base plate, or the like). In particular, it is strongly demanded to increase heat dissipating efficiency by reducing difference of temperature in the portions of the above-mentioned heat sink with fins.\nAs mentioned above, there is a problem in which a one-piece heat sink made of forged aluminum does not attain enough chat dissipating effect, because the required number of the heat dissipating fins are not formed due to the technical difficulty by forging to shorten the pitch between fins for the required level of fin-density. Furthermore, there is also a problem in which using the heavy copper base plate becomes difficult due to the heavy weight thereof even if the heat conductivity is excellent. Therefore, it is strongly demanded to increase heat dissipating efficiency of the heat sink with fins having construction that the heat dissipating fins are mechanically fixed on a surface of the aluminum base plate.\nHowever the heat sink having construction that the heat dissipating fins are mechanically fixed on a surface of the aluminum base plate has problem in heat dissipating efficiency, because there are temperature difference both in the base plate and the heat dissipating fins.\nThere is provided according to the invention a heat sink with fins comprising:\na group of heat dissipating fins which comprise combination of heat dissipating fins made of at least two kinds of metals having different heat conductivity; and\na metal base plate on which surface said group of heat dissipating fins are densely jointed."}
{"text": "The present invention relates to a transport container system, in particular for bulk goods, comprising a stackable transport container, preferably a crate, which may be of the collapsible or non-collapsible type. This type of crate is known and used in particular for transporting bulk goods such as fruit and vegetables. The term bulk goods as used in the context of the present invention shall denote a unit of goods to be transported which consists of discrete pieces of a minimum size between 0.5 cm and 1.0 cm.\nThe non-collapsible and collapsible containers of the prior art, in particular crates, for transporting fruit and vegetables are made of cardboard, wood or plastic. The special feature of collapsible transport containers is that their side walls can be moved down onto the inner bottom surface of the transport containers, which results in a volume reduction of the empty transport container. When folding the container up again, the side wall elements will be arranged perpendicular (at 90°) to the inner bottom surface and will be detachably connected to each other through various means. While the bottom surface of the transport containers is of a defined size, there are containers which have side walls of different heights, in which two or four side wall elements have the same height, to allow different transport volumes to be obtained. Furthermore, on their upper side facing away from the bottom surface, the side wall elements are provided with a profile or a means to make them stackable. In order to increase the stability of the transport containers, especially as regards their stackability, these are preferably reinforced at their corners. The maximum volumetric capacity of the prior art transport containers is defined by the size of the bottom surface and the height of the side wall elements. For a higher volumetric capacity, the transport containers must have different and higher side wall elements. This does not allow for a fast adjustment of their volumetric capacity to changing consumer demands.\nThe sizes of certain kinds of fruit and vegetables wilt vary from one harvest season to the next depending on different factors, for example during their growth period. The sizes of fruit or vegetables to be packaged are specified in EC regulations. The varying sizes of the bulk goods to be transported are thus a known problem in the transport of bulk goods such as fruit and vegetables which makes optimal filling of transport containers difficult. In order to cope with the varying demands posed by the bulk goods, the transport containers, in particular crates, are machine-produced in certain sizes which are also determined by the production line and/or by the production parameters selected. This makes it impossible to rapidly change the size—and thus the volumetric capacity—of a vast number of transport containers so as to ensure optimal filling of the containers based on the size of the bulk goods without major logistic transport problems or a time-consuming change-over of production lines and resulting high costs.\nThe above mentioned problem will crop up with the prior art transport containers especially when relatively easy-to-produce cardboard packaging for transporting bulk goods such as fruit and vegetables is replaced with returnable containers made of plastic or a material similarly suitable for this purpose which are friendlier to the environment but also more complex and costly in production. An ideal adjustment of the transport containers to the size of the bulk goods to be transported will prove especially complex and difficult in the case of the prior art returnable plastic containers. For maximum utilization of the means of transport, the transport containers can be stacked which allows a vast number of them to be transported in large containers, on loading areas, in goods wagons or similar means of transport. The bulk goods thus transported must not protrude above the upper edge of the transport containers since this would interfere with the stacking of the transport containers or otherwise damage the bulk goods. As a consequence, the volumetric capacity of the prior art transport containers cannot be fully utilized in many instances.\nThe applicant's returnable transport containers, the technical term for which is “round trip containers”, come in at least ten different designs which differ in the height of their side wall elements. The heights of the side wall elements range between 8 cm and 28 cm, with heights of 8 cm, 10 cm, 13 cm, 15 cm, 16 cm, 18 cm, 20 cm and 23 cm being preferably used. The bases of these transport containers are preferably rectangular in shape and their external measurements are preferably 600 mm×400 mm. This is approximately an integer fraction of the size of the surface area of standard Euro and U.S. pallets. However, transport containers of a different size, for example 400 mm×300 mm, are also used.\nNL 93 00 986 discloses a container having at least one bottom element and one wall element. Provided on the wall element are projections which can be made to engage in recesses provided in the bottom element to connect the wall element to the bottom element. It is furthermore disclosed in NL 93 00 986 how a circumferential one-piece frame may be placed on a container to enlarge its volume which is delimited by the wall element and the bottom element.\nDE 103 26 574 A1 discloses a transport container, in particular for transporting bulk goods such as fruit and vegetables, comprising a collapsible or non-collapsible stackable crate having a bottom element as well as four side wall elements of a pressure- and/or load-resistant structure. For increasing the volumetric capacity of the transport container, an attachment unit is provided whose shape corresponds to that of the side wall elements and which can be placed on top of the side wall elements of the crate. The attachment unit has been designed to form a closed frame which can be folded at its diagonal corners. The foldable attachment unit is preferably made of cardboard and will be disposed of after use. The side wall elements of the crate and the side wall elements of the attachment unit can be snapped into mutual engagement when the attachment unit is put on top of the crate, and can be released again when the attachment unit is taken off.\nUS 2004/0222222 A1 discloses a collapsible transport container which is adjustable in height. The transport container has a base which also constitutes the bottom surface of the transport container. The transport container furthermore includes a pair of long side walls extending opposite each other and a pair of short side walls extending opposite each other, with extension walls being provided on each of the side walls. Together with the extension walls, the short and the long side walls can be folded down onto the base to reduce the volume of the empty transport container to a minimum. The short side walls and the long side walls can be arranged so as to extend perpendicular to the bottom surface in which position they will then be locked with each other by means of locking elements provided on the short side walls. If required, long extension walls may be folded out from the long side walls and short extension walls may be folded out from the short side walls so as to form—in a first embodiment—upwards extensions each of the long and of the short side walls. Once folded out, the extension walls will be mutually locked, by means of additional locking elements disposed in the short extension walls, so as to form a frame.\nThe attachment unit can thus be taken off and disposed at the place of delivery. Once it has been emptied and cleaned, the reusable crate may for example be folded and stacked and will then be ready for future use for which no attachment units, attachment units of a different height or the same attachment units but a different amount thereof may be required. In most cases, it therefore makes more sense to store the crates separately from the attachment units. This leads to various costs, on the one hand for producing the attachment units and on the other hand for storing crates and attachment units separately, with the increased expenditure being incurred both at the place where the crates are filled and at the place where they are made. Additional costs will also be incurred at the place where the crates with the attachment units are emptied, due to the disposal of the cardboard attachment units."}
{"text": "Nutritional mineral deficiencies (also referred to as micronutrient malnutrition) affect a large proportion of the human population and are particularly widespread in developing countries. Deficiency in iron is especially common, affecting more than two billion people world-wide, and is the only nutrient deficiency that is significantly prevalent in industrialised countries. Iron deficiency is associated with a range of health problems, including impairment of work performance, increased maternal and child mortality, and poor cognitive development in children.\nIron deficiencies originate when physiological requirements are not met by mineral absorption from the diet, for example due to low iron bioavailability. Dietary iron bioavailability may be low, for example, in populations having monotonous plant-based diets with little meat. Indeed, a major cause of low iron bioavailability is related to the presence of anti-nutritional factors that are naturally present in cereals and legumes. These factors form insoluble complexes with iron and interfere with its absorption.\nFortifying food with iron is a well known approach to increase dietary iron intake, but can pose a number of difficulties. A significant problem is caused by the general incompatibility between bioavailability and stability of iron compounds. Typically, the most bioavailable iron compounds are the most reactive within the food matrix. As such, fortifying foods with iron can lead to a number of undesirable changes in properties of the food, in particular organoleptic properties of the food. For example iron can accelerate oxidation reactions, adversely altering a food's flavour, and iron can also form complexes with phenolic chromophore compounds, leading to unwanted colour changes in the food.\nFor example, ferrous (i.e. Fe(II)) sulfate, which is a reference iron compound for food fortification in humans, causes sensory changes in the food vehicle in the presence of polyphenols or high amounts of lipids. In contrast, more stable iron sources, which are typically water-insoluble (e.g. ferric (i.e. Fe(III)) pyrophosphate), have relatively low bioavailability compared to water-soluble compounds.\nA number of approaches have been taken during attempts to improve iron fortification of foods. Encapsulated ferrous sulfate has been considered, because it could provide a highly bioavailable iron source while maintaining stability through encapsulation of the formulation. However, bioavailability is highly dependent on the coating used and in many cases bioavailability of the coated iron source is reduced. Encapsulation also increases production costs. Moreover, most coatings used for encapsulation, which are lipid based, give rise to problems associated with melting during the different heat treatment stages of the manufacture of many food products.\nAlternative approaches have used iron-containing nanoparticles that have been stabilised with biopolymers (EP 1743530), or ferric EDTA (ethylenediaminetetraacetic acid), which has good bioavailability and stability (U.S. Ser. No. 10/969,434; published as US 2005/0053696), for iron fortification. However, the use of nanotechnology and EDTA in food products is meeting with increased consumer resistance. Furthermore, sodium iron EDTA is expensive and not stable in all food matrices (e.g. chicken bouillon).\nAccordingly, there remains a significant need for compositions and methods that enable fortification of foods and beverages with iron. In particular, there remains a need for compositions and methods that provide a soluble, preferably bioavailable source of iron that has minimal effect on the organoleptic properties of foods."}
{"text": "1. Field of the Invention\nThe present invention relates generally to golf equipment accessories and, more specifically, to a golf club protection device which can be inserted into the readily available individual golf club tubes which are sold and inserted into a golf bag to segregate and protect individual golf clubs. While these tubes do segregate the golf clubs one from the other they do not prevent the shafts from scraping the lip of the tube while being extracted from or inserted into during the course of play. This contact with the lip of the tube will remove the coating on graphite shafts.\nThe present invention comprising a flanged insert seats into the top opening of existing tubes by means of compression fitting or is therein permanently affixed by adhesive means and said flanged insert has an acrylic soft pile material circumferentially affixed to the interior wall of said flanged insert.\nAnother embodiment is provided having all of the properties of the preferred embodiment in addition to a club head bearing member comprised of a Y-shaped member fixedly attached to the exterior surface of the flanged insert therein providing means for supporting the club head within the V-portion and preventing any lateral movement of the club head which normally occurs during the course of play.\n2. Description of the Prior Art\nThere are other golf club shaft protection devices. Typical of these is U.S. Pat. No. 4,938,349 issued to Burns on Jul. 3, 1990.\nAnother patent was issued to Henry et al. on Jan. 4, 1994 as U.S. Pat. No. 5,275,278. Yet another U.S. Pat. No. 5,632,690 was issued to McConville on May 27, 1997 and still yet another was issued on Feb. 24, 1998 to King et al. as U.S. Pat. No. 5,720,388.\nA golf club protector for protecting the shaft of a club including a golf club housing tube having a protective interior and a protective collar provided at the upper end of the tube. Alternatively, an optional protective sleeve insertable within the tube having a protective collar attached thereto may be removably provided as a unit for a conventional tube. Where the housing tube includes a fixed collar and protective interior, the length of the tube may be predetermined in accordance with the length of a particular club shaft, or the tube may be provided with spaced apart cut markings at its lower end to facilitate shortening of the tube by the consumer.\nA golf club shaft protector is disclosed as including a hollow plastic tube of predetermined diameter and length with open upper and lower ends. The length of the hollow plastic tube substantially corresponds to the length of a golf club shaft and can be pre-selected and cut to the appropriate length. A flexible restricted throat element extends across the open upper end of the hollow plastic tube for resilient deformation upon the insertion of a golf club handle to allow passage of the golf club handle and associated golf club shaft into the hollow plastic tube. The flexible restricted throat element subsequently returns to its initial shape for close fitting circumferential support of the golf club hosel at an upper end of the golf club shaft adjacent the golf club head. The predetermined diameter of the hollow plastic tube is dimensioned to at least substantially peripherally engage the golf club handle at least adjacent the open lower end. The entire length of the golf club shaft is thus protected by the hollow plastic tube through the circumferential support of the golf club hosel by the flexible restricted throat opening at the open upper end and by the at least substantial peripheral engagement of the golf club handle at the lower open end so as to provide suspended non-engagement of the golf club shaft therebetween.\nA wrap for protecting a portion of a golf club shaft from abrasion within a golf bag. The inventive device includes a pad strip for circumferentially extending about a shaft of a golf club. A securing strip extends from the pad strip for securing the pad strip in an annular configuration about the shaft to protect the shaft from abrasion against an upper interior surface of a receiving tube of a golf bag.\nA golf club shaft protector is disclosed as including an elongated hollow plastic tube having a polygonal cross-sectional configuration with an unequal number of sides interconnected to each other by corner sections. Each of the unequal number of sides have the same predetermined length and each of the corner sections have the same predetermined angular shape. One of each of the corner sections faces one of each of the sides. Preferably, the unequal number of sides is at least seven to provide the largest possible opening with the greatest amount of rigidity for receiving the golf club shaft, including a golf club handle. The golf club shaft protector also includes a tubular element mounted adjacent the open upper end of the hollow plastic tube and includes an outer wall section, an inner wall section and flexible lip sections. The outer wall section surrounds an outer wall area of the tube adjacent the open upper end, the inner wall section surrounds an inner wall area of the tube adjacent the open upper end and the flexible lip sections extend over the open upper end of the tube for resiliently supporting a golf club shaft in centered position relative to the hollow plastic tube. The aforementioned elements facilitate the entry and removal of a golf club shaft including a golf club handle or grip, while protecting the golf club shaft against marring, scratching or other damage while retained within the hollow plastic tube.\nWhile these golf club shaft protection devices may be suitable for the purposes for which they were designed, they would not be as suitable for the purposes of the present invention, as hereinafter described.\nThe present invention discloses a device for protecting golf clubs from being scratched due to the golf clubs being inserted into a golf club protection tube which are commercially available. The present invention discloses a flanged tubular insert for placement into the top neck of the golf club protection tube. The present invention is padded with soft pile material extending circumferentially about the insert having a hole therein for inserting of the golf club shaft. Another embodiment is provided which has a Y-shaped upstanding member mounted on the top edge of the present invention within which is placed the golf club head.\nA primary object of the present invention is to provide a flanged insert which can be used in conjunction with existing golf club tubes for the protection of golf clubs.\nAnother object of the present invention is to provide a flanged insert which can be used in conjunction with existing golf club tubes to prevent the golf club shafts from contact with the lip of existing golf club tubes.\nYet another object of the present invention is to provide a golf club tube flanged insert having a lower portion of smaller diameter and an upper portion of greater diameter.\nStill yet another object of the present invention is to provide a golf club tube flanged insert having a soft pile material circumferentially affixed to the greater upper interior wall of said flanged insert.\nYet another object of the present invention is to provide an alternate embodiment of a golf club tube flanged insert having a soft pile material circumferentially affixed to the greater upper interior wall of said flanged insert and further having a club head bearing member comprised of a Y-shaped member fixedly attached to the exterior surface of the flanged insert therein providing means for supporting the club head within the V-portion and preventing any lateral movement of the club head which normally occurs during the course of play.\nAdditional objects of the present invention will appear as the description proceeds.\nThe present invention overcomes the shortcomings of the prior art by providing a golf club protection device comprising a flanged insert which seats by compression fitting or adhesively fixed into the top opening of existing golf club protection tubes and said flanged insert having an acrylic soft pile material circumferentially affixed to the interior wall surface of said flanged insert.\nAnother embodiment is provided having all of the properties of the preferred embodiment in addition to a club head bearing member comprised of a Y-shaped member fixedly attached to the exterior surface of the flanged insert therein providing means for supporting the club head within the V-portion and preventing any lateral movement of the club head which normally occurs during the course of play.\nThe foregoing and other objects and advantages will appear from the description to follow. In the description reference is made to the accompanying drawings, which form a part hereof, and in which is shown by way of illustration specific embodiments in which the invention may be practiced. These embodiments will be described in sufficient detail to enable those skilled in the art to practice the invention, and it is to be understood that other embodiments may be utilized and that structural changes may be made without departing from the scope of the invention. In the accompanying drawings, like reference characters designate the same or similar parts throughout the several views.\nThe following detailed description is, therefore, not to be taken in a limiting sense, and the scope of the present invention is best defined by the appended claims."}
{"text": "The term “biocide” or “biocidal” refers to a chemical substance capable of killing living organisms, usually in a selective way. Biocides are commonly used in medicine, agriculture, forestry, and in industry where they prevent the fouling of water and oil pipelines. Some substances used as biocides are also employed as anti-fouling agents or disinfectants under other circumstances. The term “biological agent” refers to any living organisms or the materials derived from them (including, but not limited to, bacteria, viruses, fungi, and toxins) that cause disease in or harm to humans that could be used in biological warfare. Harmful biological agents include anthrax (Bacillus anthracis) spores. The term “agent” refers to biological agents.\nPast agent defeat projects have used hydrogen chloride as a chemical neutralizer, e.g., Lockheed Martin's “Agent Defeat Warhead Device”: “These propellants produce gas-phase water and hydrogen chloride that combine to form very reactive hot hydrochloric acid; and as the reactants cool the cooled hydrochloric acid remains in the bunker and may act to continue neutralization of bunker contents for many days . . . . Thus, incendiary agents based on standard composite rocket propellant technology are logical choices for the application described herein.” [Jones, J. W. U.S. Pat. No. 6,382,105 (2002)]. It has long been known that HCl has among the poorest bactericidal activities of common acid species. [Paul, T. et al. Biochemische Zeitschrift 1911, 29, 202].\nOther agent defeat concepts have employed elemental chlorine (a more efficient biocide than HCl) [http://www.globalsecurity.org/military/systems/munitions/adw.htm]. Such formulations so far have produced only low levels of Cl2. One drawback of Cl2 deployment is residual toxic chloro-organics such as chloramines, etc.\nIn contrast to HCl, hydrogen fluoride (HF) is an efficient antibacterial, antimicrobial biocide: 200 ppm aq. HF destroys Pseudomonas aeruginosa and Saccharomyces diastaticus in 5 min [Bessems, E.; Junger, R. GB Patent 1584845 (1981)]; and 1-100 ppm aq. HF sterilizes water lines [Tatsuno, T. et al. U.S. Pat. No. 5,147,605 (1992)].\nNovel energetic materials have been sought whose deployment would produce biocidal reaction products following explosive events that destroy a structure or containers involving harmful chemical or biological agents, thus neutralizing the agents post-blast. Based on the attractiveness of hydrogen fluoride as a superior agent defeat by-product—from its known biocidal activity—a class of energetic ingredient that appeared particularly promising for incorporation into explosive formulations for the specific application of agent defeat weapons was the class of energetic difluoramines.\nEnergetic materials have been sought whose deployment would produce biocidal reaction products following explosive events that destroy a structure or containers involving harmful biological agents, thus neutralizing the agents post-blast. Thus, there exists a need in the art for such an explosive that is has high biocidal activity during and/or after detonation, and is an easily synthesized, powerful and efficient chemical neutralizer.\nIt is to be understood that the foregoing general description and the following detailed description are exemplary and explanatory only and are not to be viewed as being restrictive of the invention, as claimed. Further advantages of this invention will be apparent after a review of the following detailed description of the disclosed embodiments, which are illustrated schematically in the accompanying drawings and in the appended claims."}
{"text": "As digital apparatuses provide diverse functions nowadays, digital apparatuses having a built-in camera are developed and used. Recently, there emerges a digital apparatus that automatically moves a lens of a camera built in the digital apparatus using a predetermined lens driving motor, thereby allowing a user of the digital apparatus to use the camera more conveniently.\nThe lens driving motor should be provided in a small size with consideration of miniaturization of the digital apparatus. Also, an inner portion of the lens driving motor, particularly, an electrical line through which a current flows should be stably connected, and reliability in accuracy for a movement amount of a lens should be excellent considering that the lens driving motor is an optical apparatus. Japanese Patent Publication No. 2004-280031 proposes a related art lens driving motor capable of the above-described purposed to some extent.\nThe Japanese Patent discloses a lens driving motor including a coil and a carrier coupled together, a lens coupled to the carrier, and a magnet installed on an outer side of the coil.\nThis related art lens driving motor operates in such a way that the carrier and the lens move in an upward direction using electromagnetic force generated between the coil and the magnet when a current flows through the coil. Also, one side of a spring is fixed to the carrier, and the other side of the spring is clamped by a yoke and caps of upper and lower portions so that a stop position of the carrier is indicated.\nMeanwhile, according to the Japanese Patent, lead lines of a coil through which a current is applied to the coil are soldered at an upper spring and a lower spring, respectively, to receive external power via the upper spring and the lower spring.\nHowever, the lead lines of the coil should be soldered at the upper and lower springs, respectively, during a process of assembling respective parts because of the structure of the related art lens driving motor. Since a lens driving apparatus has a very small volume of about 1.2-1.5 cm3, a soldering operation should be performed in an inside of a narrow space, which reduces workability.\nAlso, due to a yoke shape where cross-sections are bent, the lead line of the coil should be inserted into a gap between the yoke and the carrier, extracted to an upper side, and then soldered at the upper spring when the lead line is soldered at the upper spring. These processes are more difficult to perform.\nAlso, since one end of the lead line of the coil is connected to the coil, of course, and move together with the coil, but the other end of the coil is fixed to a portion (e.g., the yoke) that moves independently of the coil, friction caused by contact and an external load are applied to the lead line. These problems generate deterioration of the lead line and reduce reliability of a product."}
{"text": "A wristwatch, as a timing instrument, is always an indispensable commodity for people in their daily work and life. A wristwatch usually comprises: a time indicating component which is the kernel component of the wristwatch and is used for indicating time; a watch case for housing and protecting the time indicating component; a watch band for fixing the wristwatch to people's wrist and for adorning purpose. With the development of science and technology, people integrate other devices into the wristwatch to make full use of the wristwatch as a portable commodity. For example, in a patent application numbered “CN01 104642.2” and titled “a telephone wristwatch with a SIM card”, a wristwatch combined with communication technology is disclosed. But, there is still no portable wristwatch combined with computer technology up to now, despite of the wide spreading and continuous development of computer technology. On the other hand, to meet the desires of data exchange, various types of mobile memory devices are present in the market, such as mobile hard disks, portable memory devices with a USB connectors, etc. Data exchange can be realized between computers by utilizing these products. For example, data from one computer is downloaded and stored to a portable memory device with a USB connector quickly, and then the data in the memory device is transmitted to another computer through the USB connector. Generally, the data stored in these products with Flash Memory can be downloaded rapidly and stored for a long time, for some instance, more than ten years. However, the mobile hard disk is not portable, and the existing portable memory device with a USB connector has to be equipped with a long cable, which connects with the USB connector of the product on its one end, and with the USB port of the computer on the other end. It often causes a lot of inconvenience since a consumer has to always carry the cable with him, and data transmission and storage becomes impossible if the cable is lost. Therefore, it is desirable to develop a portable memory device integrated with daily commodity like a wristwatch, so that the USB product needs not to be carried additionally."}
{"text": "The present disclosure relates to producing resources from a subterranean zone.\nOften, an injection treatment will be applied to a well prior to putting the well on production or at some point during production. Some example injection treatments include acidizing or solvent injection to remove near wellbore damage, steam injection to mobilize resources in a formation, and water or polymer-laden fluid sweeping to pressurize and sweep resources in a reservoir to a desired location. There are other types of injection treatments. Because it is costly and time consuming to run different well strings into and out of a wellbore, the injection treatments are performed with the production string when practicable."}
{"text": "Currently available and well known plastic lined piping products comprise a family of pipes, fittings, and valves especially designed for handling corrosive or high purity liquids. Such products generally comprise steel lined with a polymeric material, e.g., polyvinylidene chloride, polypropylene, polyvinylidene fluoride, and polytetrafluoroethylene. Such products enjoy both the structural integrity of steel and the high chemical resistance characteristic of the selected polymeric liner. The external surface of such products often bears a thin coating of paint to protect the steel from corrosive substances found in the operating environment.\nEach pipe, fitting, and valve will contain a flange at or substantially near each end thereof. Adjacent pipes, fittings, and valves, within a given pipeline, may be joined one to the other by the fastening together of such flanges, e.g., by bolting. Such bolts may be coated with a polymer, e.g., polytetrafluoroethylene, to inhibit the corrosion thereof.\nTo secure the polymeric liner within a given pipe, fitting, or valve, such liner is provided with a length greater than that of the bore through which it extends. Upon the application of heat and pressure, the liner is flared over the pipe end, which in some cases will be the anterior flange face. In such cases, the flaring will preferably not cause an extension of the liner to the circumferential edge of the flange. Thus, when adjacent flanges are secured together, e.g., by bolting, a gap remains between adjacent opposing anterior flange faces equal to at least the thicknesses of the flared polymeric liners. The presence of this gap precludes the conductivity of electricity between adjacent components when non-conductive bolts, e.g. polytetrafluoroethylene coated bolts, are used.\nMethods of establishing electrical conductivity between adjacent components are known in the art. For instance, external lock tooth washers may be compressed between the bolt head and posterior flange face of a first component, and between the nut and flange of the adjacent component. Such washers bear teeth which serve to remove the paint from the outer surface of the posterior flange faces and the non-conductive coating from the bolts, thereby providing points of conductivity. While providing improved conductivity along a given pipeline, external lock tooth washers present undesirable complications. First, the installation of such washers or the replacement of a corroded washer after the assembly of the pipeline is impossible without dismantling at least the effected portion of the pipeline. This often involves clearing the pipeline of chemical substances contained therein, removing the connecting bolts, installing the washer(s), and retorquing the bolts. Second, the presence of external lock tooth washers causes torquing problems. Such washers cause an increase in friction between the nut and the posterior flange face. This increased friction is added torque, causing an inaccurate torque reading by the torque wrench. Thus, such washers may lead to insufficiently tight bolts, which may in turn lead to leaks at the pipe joint.\nIt is further known to provide conductivity clamps on adjacent components which bite into the paint and provide points of conductivity. A conductive wire connects the clamps located on adjacent pipes, providing a conductive path therebetween. Such clamps tend to be relatively expensive.\nThose in the industry would find great advantage in a means for providing electrical contact between the opposing anterior faces of adjacent metal pipe flanges which is both cost effective and practical. Such means should be easy replaceable and should not result in torquing problems. The loss or corrosion of such means should likewise be easily detectable. Such means should preferably permit the establishment of electrical conductivity between relatively widely spaced anterior flange faces, e.g., anterior flange faces spaced approximately one inch apart."}
{"text": "1. Field of the Invention\nThis invention relates to a new, compact, and improved negative ion source and injector system, including preaccelerator and beam transport innovations. More particularly, it relates to methods and apparatus for efficiently excluding electrons from the negative ion beam and creating a high-angle convergent beam for injection into other devices including accelerator stages.\n2. Description of the Background\nEnergetic neutral beams are difficult to create but are necessary for many applications. For example, charged particles cannot be injected through the large magnetic fields surrounding many high energy accelerators, for example the proton synchrotron, and other devices because the charged particles would be deflected from their path by the magnetic field. None-the-less, particles must be charged so that a voltage gradient can be used to extract them from the plasma. This means that particles used in high energy devices are likely to go through several stages prior to injection in the device. They are initially electrically charged so they can be extracted from the plasma gas that contains the negative ion. The charged particles are accelerated through many stages, then the beam has its charge changed just prior to injection into the high energy device.\nIt is desirable that as many electrons as possible are eliminated from the beam immediately upon exit from the ion source. The undesired high energy electrons, found in negative ion beams, are currently separated from a negative hydrogen ion beam by using an ExB electron extractor which diverts them to collector electrodes. Because of the magnetic field surrounding them, permanent magnets could not be placed near the exit aperture of the ion source without compromising the efficiency of the source. Placing the magnets downstream requires a considerable amount of power and is a substantial power drain on a system.\nOne approach for the production of negative ions and suppression of accompanying electrons is disclosed in US. Pat. No. 4,486,665. The approach described therein filters high-energy electrons by dividing the ion source chamber into two zones, an ionizing zone and an extraction zone. Excess electrons are suppressed by placing a plasma grid adjacent to the extraction zone of the ionization vessel that is positively-biased with respect to the anode. The positively biased grid suppresses electrons that would be dragged along with the positive ions which originate in the ionizing zone and pass through the magnetic filter. Additionally, electrons are suppressed from the output by providing a magnetic field, aligned with respect to the electric field of the extractor, to provide ExB drift to the electrons. There is some difficulty with that system, however. The magnets described in that patent are located interior to the ion source. Those magnets interfere with negative ion extraction, reducing the number of negative ions extracted by approximately a factor of 2. In addition, fewer electrons are successfully diverted from the negative ion beam than with the present invention. Further, when a radio frequency source is used to generate the negative ions, the magnet field resulting from the magnets as used in US. Pat. No. 4,486,665 interfere with the production of negative ions within the ion source volume. In order to create a magnetic field to deflect electrons from the extracted beam without interfering with ion production and extraction for the invention described in that patent, the magnets would have be located far enough down the beam transport line that the magnetic field would not penetrate the ion source vessel. This would lengthen the beam transport dimensions and interfere with the geometry of the electrodes used to focus the beam. If the magnetic field does enter the ion source vessel, negative ion production becomes less efficient. Lengthening the beam transport dimensions requires added power consumption and control circuitry.\nIt is important to keep the gas pressure low in the accelerator column. If gas pressure rises too high, collisions between negative ions and contaminants will cause electrons to be stripped from the ion prematurely.\nThere are several examples of devices that require energetic neutral beams created from negative ion particles. For example, large magnetic-confinement fusion energy devices such as tokamac accelerators, require 1 MeV neutral beams for injection. In order for the extracted beam to be injected into the cyclotron storage and acceleration ring without being defocused when crossing the electromagnetic fields, the beam charge must be changed and negative ions are more easily changed than positive ions. High energy cyclotron accelerators also require charge changed beams in order for the beam to be extracted from the cyclotron, booster ring. (L. W. Alvarez, Rev. Sci. Instr. 22(1951)705). Another application of the invention is for energetic particle beams that must be neutral in charge to transit large distances in space. The superconductor supercollider (SSC) is similarly in need of a high energy particle beam that could be created from a high quality negative ion source."}
{"text": "Lean premixed combustion systems have been deployed on land based gas turbine engines to reduce emissions, such as NOx and CO. These systems have been successful and, in some cases, produce emission levels that are at the lower limits of measurement capabilities, approximately 1 to 3 parts per million (ppm) of NOx and CO. Although these systems are a great benefit from a standpoint of emission production, the operational envelope of the systems is substantially reduced when compared to more conventional combustion systems. As a consequence, the control of fuel conditions, distribution and injection into the combustion zones has become a critical operating parameter and requires frequent adjustment, when ambient atmospheric conditions, such as temperature, humidity and pressure, change. The re-adjustment of the combustion fuel conditions, distribution and injection is termed tuning.\nControlled operation of a combustion system generally employs a manual setting of the operational parameters of a combustor at an average operational condition. These settings are satisfactory at the time of the setup, but conditions may change and cause an unacceptable operation in a matter of hours or days. Other approaches use a formula to predict emissions based on gas turbine operating parameters and select a set point for fuel distribution and/or overall machine fuel/air ratio, without modifying other parameters, such as fuel gas temperature. These approaches do not allow for timely variation, do not take advantage of actual dynamics and emission data or do not modify fuel distribution, fuel temperature and/or other turbine operating parameters.\nAnother variable that impacts the lean premixed combustion system is fuel composition. Sufficient variation in fuel composition will cause a change in the heat release of the lean premixed combustion system. Such change may lead to emissions excursions, unstable combustion processes, or even blow out of the combustion system.\nMis-operation of the combustion system manifests itself in augmented pressure pulsations or an increase in combustion dynamics. Pulsations can have sufficient force to destroy the combustion system and dramatically reduce the life of combustion hardware. Additionally, improper tuning of the combustion system can lead to emission excursions and violate emission permits. Therefore, a means to maintain the stability of the lean premixed combustion systems, on a regular or periodic basis, within the proper operating envelope, is of great value and interest to the industry. Additionally, a system that operates by utilizing near real-time data, taken from the turbine sensors, would have significant value to coordinate modulation of fuel distribution, fuel gas inlet temperature and/or overall machine fuel/air ratio."}
{"text": "(a) Field of the Invention\nThis invention relates to film winding devices adapted to slide projectors of a type using strip films.\n(B) Description of Prior Art\nThere is already a type of slide projector of a type of using strip films wherein a strip film wound in the form of a roll is placed on the upper portion of a projector body and is wound into a housing part provided in the lower portion of the projector body by means of a sprocket means through a film gate. However, in the case of such type, when a once projected film is reprojected as it is, the projecting order will be reversed. Therefore, in practice, it is usual to rewind the projected film. As a result, this enforces a trouble more than is necessary on the operator in practice."}
{"text": "Aqueous cleaning processes are a mainstay of both domestic and industrial textile fabric washing. On the assumption that the desired level of cleaning is achieved, the efficacy of such processes is usually characterised by their levels of consumption of energy, water and detergent. In general, the lower the requirements with regard to these three components, the more efficient the washing process is deemed. The downstream effect of reduced water and detergent consumption is also significant, as this minimises the need for disposal of aqueous effluent, which is both extremely costly and detrimental to the environment.\nSuch washing processes involve aqueous submersion of fabrics followed by soil removal, aqueous soil suspension, and water rinsing. In general, within practical limits, the higher the level of energy (or temperature), water and detergent which is used, the better the cleaning. The key issue, however, concerns water consumption, as this sets the energy requirements (in order to heat the wash water), and the detergent dosage (to achieve the desired detergent concentration). In addition, the water usage level defines the mechanical action of the process on the fabric, which is another important performance parameter; this is the agitation of the cloth surface during washing, which plays a key role in releasing embedded soil. In aqueous processes, such mechanical action is provided by the water usage level in combination with the drum design for any particular washing machine. In general terms, it is found that the higher the water level in the drum, the better the mechanical action. Hence, there is a dichotomy created by the desire to improve overall process efficiency (i.e. reduce energy, water and detergent consumption), and the need for efficient mechanical action in the wash. For domestic washing in particular there are defined wash performance standards specifically designed to discourage the use of such higher levels in practice, in addition to the obvious cost penalties which are associated with such usage.\nCurrent efficient domestic washing machines have made significant strides towards minimising their consumptions of energy, water and detergent. EU Directive 92/75/CEE sets a standard which defines washing machine energy consumption in kWh/cycle (cotton setting at 60° C.), such that an efficient domestic washing machine will typically consume<0.19 kWh/kg of washload in order to obtain an ‘A’ rating. If water consumption is also considered, then ‘A’ rated machines use <9.7 liters/kg of washload, whilst the most efficient modern machines are now capable of using even less water—e.g. model number F1480FD6 manufactured by LG (see www.lg.com). This machine typically uses 63 liters for a 9 kg washload, i.e. 7 liters/kg.\nDetergent dosage is then driven by manufacturer recommendations but, again, in the domestic market, for a concentrated liquid formulation, a figure of 35 ml (or 37 g) for a 4-6 kg washload in soft and medium hardness water, increasing to 52 ml (or 55 g) for a 6-8 kg washload (or in hard water or for very dirty items) is typical (see, for example, Unilever pack dosage instructions for Persil® Small & Mighty). Hence, for a 4-6 kg washload in soft/medium water hardness, this equates to a detergent dosage of 7.4-9.2 g/kg whilst, for a 6-8 kg washload (or in hard water or for very dirty items), the range is 6.9-9.2 g/kg.\nEnergy, water and detergent consumptions in the industrial washing process (washer-extractors) are considerably different, however, and usages of energy and water are less constrained in such environments, since these are principal factors in reducing cycle time—which is, of course, more of a consideration than in the domestic scenario. There is a similar pressure on detergent levels, however, but this is mostly due to a desire to reduce cost.\nThus, it can be taken from the above discussion that the performance levels which set the highest standard for an efficient fabric washing process are an energy consumption of <0.19 kWh/kg, a water usage of approximately 7 liters/kg, and a detergent dosage of approximately 8 g/kg. However, as already mentioned, it is becoming increasingly difficult to reduce the water (and, hence, energy and detergent) levels in a purely aqueous process, due to the minimum requirement to wet the fabric thoroughly, the need to provide sufficient excess water to suspend the soil removed in an aqueous liquor and, finally, the need to rinse the fabric.\nHeating of the wash water is then the principal use of energy, and a minimum level of detergent becomes necessary in order for an effective concentration to be reached at the operating wash temperature. If a means to improve mechanical action could be achieved without increasing the water level used, then the aqueous wash process could become significantly more efficient (i.e. yield further reductions in energy, water and detergent consumption). It should be noted that mechanical action itself has a direct effect on the detergent level, since the greater the level of soil removal which is achieved through physical force, the less that is required of the detergent chemistry. However, increasing the mechanical action in a purely aqueous washing process has certain associated drawbacks. Fabric creasing readily occurs in such processes, and this acts to concentrate the stresses from mechanical action at each crease, resulting in localised fabric damage. Prevention of such fabric damage (i.e. fabric care) is of primary concern to the domestic consumer and the industrial user.\nVarious different approaches to the development of new cleaning technologies have been reported in the prior art, including methods which rely on electrolytic cleaning or plasma cleaning, in addition to approaches which are based on ozone technology, ultrasonic technology or steam technology. Thus, for example, WO-A-2009/021919 teaches a fabric cleaning and disinfection process which utilises UV-produced ozone along with plasma. An alternative technology involves cold water washing in the presence of specified enzymes, whilst a further approach which is particularly favoured relies on air-wash technology and, for example, is disclosed in US-A-2009/0090138. In addition, various carbon dioxide cleaning technologies have been developed, such as the methods using ester additives and dense phase gas treatments which are described in U.S. Pat. No. 7,481,893 and US-A-2008/0223406, although such methods generally find greater applicability in the field of dry cleaning. Many of these technologies are, however, technically complex and not readily suited to domestic applications, in particular.\nIn the light of the challenges which are associated with aqueous washing processes, the present inventors have previously devised a new approach to the problem, which is technologically straightforward, and yet still allows the deficiencies demonstrated by the methods of the prior art to be overcome. The method which is provided eliminates the requirement for the use of large volumes of water, but is still capable of providing an efficient means of cleaning and stain removal, whilst also yielding economic and environmental benefits.\nThus, in WO-A-2007/128962 there is disclosed a method and formulation for cleaning a soiled substrate, the method comprising the treatment of the moistened substrate with a formulation comprising a multiplicity of polymeric particles, wherein the formulation is free of organic solvents. Preferably, the substrate is wetted so as to achieve a substrate to water ratio of between 1:0.1 to 1:5 w/w, and optionally, the formulation additionally comprises at least one cleaning material, which typically comprises a surfactant, which most preferably has detergent properties. In preferred embodiments, the substrate comprises a textile fibre and the polymeric particles may, for example, comprise particles of polyamides, polyesters, polyalkenes, polyurethanes or their copolymers, but are most preferably in the form of nylon chips.\nThe use of this cleaning method, however, presents a requirement for the cleaning chips or beads to be efficiently separated from the cleaned substrate at the conclusion of the cleaning operation, and this issue was initially addressed in WO-A-2010/094959, which provides a novel design of cleaning apparatus requiring the use of two internal drums capable of independent rotation, and which finds application in both industrial and domestic cleaning processes.\nWith a view to providing a simpler, more economical means for addressing the problem of efficient separation of the cleaning media from the substrate at the conclusion of the cleaning process, however, a further apparatus is disclosed in co-pending PCT Patent Application No. PCT/GB2010/051960. The apparatus of PCT Patent Application No. PCT/GB2010/051960, which finds application in both industrial and domestic cleaning processes, comprises a perforated drum and a removable outer drum skin which is adapted to prevent the ingress or egress of fluids and solid particulate matter from the interior of the drum. The cleaning method requires attachment of the outer skin to the drum during a first wash cycle, after which the skin is removed prior to operating a second wash cycle, following which the cleaned substrate is removed from the drum.\nThe apparatus and method of PCT Patent Application No. PCT/GB2010/051960 is found to be extremely effective in successfully cleaning substrates, but the requirement for the attachment and removal of the outer skin detracts from the overall efficiency of the process and the present inventors have, therefore, sought to address this aspect of the cleaning operation and to provide a process wherein this procedural step is no longer necessary. Thus, by providing for continuous circulation of the cleaning chips during the cleaning process, it has been found possible to dispense with the requirement for the provision of an outer skin."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a magnetic random access memory (MRAM), and more particularly to a memory cell architecture in an MRAM that includes magnetic memory cells, each of which is formed using an element that stores “0”/“1” data by tunneling magnetoresistive effect.\n2. Description of the Related Art\nIn recent years, a variety of memories, which store data based on novel principles, have been proposed. Of these, there is known an MRAM with nonvolatility and high operation speed, wherein magnetic memory cells, each of which is formed using a magnetic tunnel junction (MTJ) element that stores “0”/“1” data by tunneling magnetoresistive (TMR) effect, are arranged in a matrix.\nFIG. 14 schematically shows a cross-sectional structure of an MTJ element 70 that is used in a conventional MRAM.\nFIGS. 15A and 15B illustrate two states of the directions of spin in two magnetic layers 71 and 72 of the MTJ element 70 shown in FIG. 14.\nThe MTJ element 70 is configured such that one non-magnetic layer (tunneling barrier film) 73 is interposed between two magnetic layers 71 and 72. The MTJ element 70 stores “0”/“1” data, depending on whether the directions of spin of the two magnetic layers 71 and 72 are parallel, as shown in FIG. 15A, or antiparallel, as shown in FIG. 15B.\nNormally, an antiferromagnetic layer 74 is disposed on one of the two magnetic layers 71 and 72. When the layer 74 is disposed on the layer 72, the antiferromagnetic layer 74 fixes the direction of spin of the magnetic layer 72. Thus, data can easily be rewritten by changing only the direction of spin of the other magnetic layer 71. The variable-spin side magnetic layer 71 is referred to as a “free layer”, and the fixed-spin side magnetic layer 72 is as a “fixed layer” (or “pinned layer”).\nAs is shown in FIG. 15A, when the directions of spin (indicated by arrows) in the two magnetic layers 71 and 72 are parallel (the same), the tunnel resistance of the tunneling barrier film 73 that is sandwiched between the two magnetic layers 71 and 72 decreases to a minimum (tunneling current increases to maximum).\nAs is shown in FIG. 15B, when the directions of spin in the two magnetic layers 71 and 72 are antiparallel, the tunnel resistance of the tunneling barrier film 73 that is sandwiched between the two magnetic layers 71 and 72 increases to a maximum (tunneling current decreases to minimum).\nFIG. 16 schematically shows an example of a plan-view layout of a memory cell array of an MRAM that incorporates conventional memory cells. This example illustrates an architecture in a data write mode.\nA plurality of write word lines WWL and a plurality of bit lines BL are arranged perpendicular to each other. At intersections of these lines, memory cells each comprising an MTJ element are disposed. Each MTJ element has rectangular shape with a longitudinal axis extending along the write word line WWL, and with a transverse axis extending along the bit line BL. The direction of spin, which is parallel to the longitudinal axis, is given to the MTJ element. In the MRAM, two states with different resistance values of the MTJ element are associated with a “1” data storage state (“1” state) and a “0” data storage data (“0” state), respectively.\nFIG. 17 is a cross-sectional view, taken along line 15—15 in FIG. 16, showing an example of the structure of, in particular, one memory cell in a cross section perpendicular to the write word line WWL.\nFIG. 18 is a cross-sectional view, taken along line 16—16 in FIG. 16, showing an example of the structure of the memory cell in a cross section perpendicular to the bit line BL.\nIn FIGS. 17 and 18, reference numeral 10 denotes a semiconductor substrate (e.g. P-type Si substrate); 11 a shallow-trench device isolation region (STI); 12 a gate oxide film; 13 an impurity diffusion layer (N+) that functions as a drain region or a source region of a read-out cell select transistor Tr (NMOSFET); 14 a gate electrode (GC); 15 a first metal wiring layer (M1); 16 a second metal wiring layer (M2); 17 an MTJ connection wire formed of a third metal wiring layer (M3); 18 a conductive contact for electrically connecting the first metal wiring 15 to the diffusion layer 13; 19 a conductive contact for electrically connecting the second metal wiring layer 16 to the first metal wiring layer 15; 20 a conductive contact for electrically connecting the third metal wiring layer 17 to the second metal wiring layer 16; 70 an MTJ element; 22 a fourth wiring layer (M4); 23 a conductive contact for electrically connecting the fourth metal wiring layer 22 to the MTJ element 70; and 24 an interlayer insulation film.\nIn the Figures, the uses of the respective wiring layers are defined as follows: (BL) is a bit line for write/read, (WWL) is a write word line, (SL) is a source line, and (RWL) is a read-out word line. The source line (SL) is connected to a ground potential.\nNow referring to FIG. 14 to FIG. 18, the operational principle of data write to the prior-art MTJ element 70 is described.\nData write to the MTJ is performed in the following manner. As is shown in FIG. 16, write currents in directions, for example, indicated by arrows, are let to flow in the write word line WWL and bit line BL. Using a composite field of magnetic fields Hy and Hx generated by these currents, the direction of spin of the free layer 71 can be set to be parallel or antiparallel, relative to the pinned layer 72. Thereby, data is written.\nFor example, when data is written to the MTJ element 70 shown in FIG. 16, a current in a first direction or in a second direction, which is opposite to the first direction, is supplied to the bit line BL in accordance with write data, thereby generating a magnetic field Hx. A current in a fixed direction is supplied to the write word line WWL, thereby generating a magnetic field Hy. Using a composite field produced by the magnetic fields Hx and Hy, data is written. In this case, if the current in the first direction is supplied to the bit line BL, the directions of spin in the MTJ element 70 become parallel. If the current in the second direction is supplied, the directions of spin become antiparallel. FIG. 16 illustrates the case where the direction of spin in the free layer 71 is made parallel to the direction of spin in the pinned layer 72 by the composite field.\nWhen data is read out of the MTJ element 70, the read-out word line RWL shown in FIGS. 17 and 18 are activated to turn on the transistor Tr that is the switching device connected to the selected MTJ element 70. Thus, a current path is formed and a current is let to flow from the selected bit line BL to ground potential. As a result, a current corresponding to the resistance value of the selected MTJ element 70 flows only through the MTJ element 70. By detecting the current value, the data can be read out.\nReferring to FIGS. 19 and 20, how the direction of spin in the MTJ element 70 is selected by the direction of applied field is described in brief.\nFIG. 19 shows variation characteristics (MTJ curve) of resistance value due to reversal of applied field in the MTJ element 70.\nFIG. 20 shows an asteroid curve of the MTJ element 70.\nAs indicated by the MTJ curve in FIG. 19, when a magnetic field Hx is applied in an easy-axis direction of the MTJ element, the resistance value (magnetoresistance (MR) ratio) of the MTJ element 70 changes by, e.g. about 17%. The ordinate in FIG. 17 expresses the resistance value of the MTJ element 70 as a change ratio (i.e. resistance ratio between pre-change and post-change). The MR ratio varies depending on the properties of the magnetic layers of the MTJ element 70. At present, an MTJ element with an MR ratio of about 50% is obtained. A composite field of an easy-axis field Hx and a hard-axis field Hy is applied to the MTJ element 70.\nAs shown by solid lines and broken lines in FIG. 19, the magnitude of the easy-axis field Hx, which is necessary for changing the resistance value of the MTJ element 70, varies depending on the magnitude of the hard-axis field Hy. The broken lines indicate MTJ curves in cases where the hard-axis field Hy is greater than in the case of the solid lines. Making use of this phenomenon, data can be written to only the MTJ element 70 in the arrayed memory cells, which is disposed at the intersection of the selected write word line WWL and selected bit line BL.\nAs is shown in FIG. 20, if the magnitude of the composite field of the easy-axis field Hx and hard-axis field Hy falls outside the asteroid curve (e.g. locations indicated by black circular marks), the direction of spin in the magnetic layer of the MTJ element 70 can be reversed.\nOn the other hand, if the magnitude of the composite field of the easy-axis field Hx and hard-axis field Hy falls inside the asteroid curve (e.g. locations indicated by white circular marks), the direction of spin in the magnetic layer of the MTJ element 70 cannot be reversed.\nAccordingly, the data write to the MTJ element 70 can be controlled by varying the magnitude of the composite field of the easy-axis field Hx and hard-axis field Hy and changing the position of the magnitude of the composite field in the Hx-Hy plane.\nIn the above-described prior-art cell architecture shown in FIG. 15, however, the MTJ element 70 is stacked via many metal layers that are provided above the read-out cell select transistor Tr. This complex stacked structure requires many conductor layers and interlayer insulation film, including eight metal wiring layers 18, 15, 19, 16, 20, 17, 23 and 22 and one MTJ element 70. Consequently, a great number of fabrication steps are needed, and it is difficult to provide an MRAM at low cost."}
{"text": "    Patent document 1: JP-2003-198713 A (corresponding to US-2003/0114202)    Patent document 2: JP-2003-218996 A\nA cellular phone compliant to a Bluetooth (registered trademark) communication function becomes widely used. In connection with it, an in-vehicle handsfree apparatus for establishing a Bluetooth communication link with the above cellular phone has been offered (for example, Patent documents 1, 2).\nIn such an in-vehicle handsfree apparatus, when the cellular phone receives an incoming signal via a communication network while establishing the Bluetooth communication link between the in-vehicle handsfree apparatus and the cellular phone, the cellular phone sends an incoming notice signal or call to the handsfree apparatus. Upon receiving it, the handsfree apparatus executes the handsfree incoming control to give a notice for indicating that the cellular phone has received the incoming call while prohibiting the cellular phone from outputting an incoming ringer tone.\nIn contrast, an in-vehicle navigation apparatus having a handsfree function connectable with multiple communication terminals is offered. A cellular phone may receive an incoming signal from a communication network in a state where the in-vehicle navigation apparatus cannot execute a handsfree incoming control relative to the incoming signal. In such a case, the navigation apparatus has a difficulty in executing a handsfree incoming control. The user needs to operate the cellular phone, instead of operating the navigation apparatus, to answer the incoming signal. Herein, it is not a problem if the navigation apparatus is operated while the vehicle is parked or stopped; however, if the vehicle is running, it becomes a cause to interfere the safety of driving the vehicle. This is not desirable."}
{"text": "The present invention relates generally to beads, methods of making beads, and methods of using beads to remove metal contaminants dissolved in aqueous solutions. More particularly, the beads preferably include peat moss and a binder and the peat moss is capable of sorbing dissolved metal ions."}
{"text": "1. Field of the Invention\nThe invention relates to a decoding circuit and a decoding method thereof, and more particularly, to a decoding circuit of a DMX512 signal format and a decoding method thereof.\n2. Description of Related Art\nDMX512 protocol standard is a data transmission standard initially proposed by United States of Institutes for Theater Technology (USITT) in 1986, and is a commonly used control protocol in the field of entertainment lighting. The DMX512 communication technology applies an asynchronous communication format, where each lighting adjustment data is composed of 11 bits, including one bit representing a start time, 8 bits representing a data byte and two stopping bits, and 512 lighting adjustment data can be transmitted every time. Currently, stage light controls mostly use the DMX512 protocol standard to perform the stage light control.\nIn the DMX512 protocol standard, every transmitting data includes a break signal and a plurality of data signals, where the break signal is usually a logic low level for 88 micro-seconds(μs). The break signal is a start signal of data, and a chip can determine a start of the signal by detecting the break signal. The break signal is followed by a plurality of time slots, and there is also a mark between every neighboring time slots as an interval, where the mark is a logic high level. When there is no DMX data packet generated, the signal is kept at the logic high level, and this is a mark after the last time slot. Therefore, there is usually a mark signal of the logic high level present before the break signal.\nNowadays, the driver circuit of light emitting diodes (LED) may also apply the DMX 512 signal format to transmit the driving data. In order to speed up the data transmission, transmission ways of 2 frequency multiplications (or two times the frequency of an standard DMX 512 protocol) or 4 frequency multiplications (or four times the frequency of an standard DMX 512 protocol) of the standard DMX512 signal format transmission rate are applied, mainly to shorten the standard time specification in the DMX512 signal format to ½ or ¼. Since the driver circuit of LED usually transfers the DMX512 signals in a serially connected fashion, if there are too many of the serially connected chips, the signal is easily made distorted, and the chip(s) serially connected thereafter thereby cannot correctly decode the data."}
{"text": "This relates generally to electronic devices and, more particularly, to electronic devices with wireless communications circuitry.\nElectronic devices often include wireless communications circuitry. For example, cellular telephones, computers, and other devices often contain antennas and wireless transceivers for supporting wireless communications.\nIt may be desirable to support wireless communications in millimeter wave communications bands. Millimeter wave communications, which are sometimes referred to as extremely high frequency (EHF) communications, involve communications at frequencies of about 10-400 GHz. Operation at these frequencies may support high bandwidths, but may raise significant challenges. For example, millimeter wave communications are often line-of-sight communications and can be characterized by substantial attenuation during signal propagation.\nIt would therefore be desirable to be able to provide electronic devices with improved wireless communications circuitry such as communications circuitry that supports millimeter wave communications."}
{"text": "1. Field of the Invention\nThe present invention relates to a spin valve sensor with an improved antiparallel (AP) pinned layer and more particularly to an AP pinned layer that has reduced current shunting and lower coercivity.\n2. Description of the Related Art\nThe heart of a computer is an assembly that is referred to as a magnetic disk drive. The magnetic disk drive includes a rotating magnetic disk, write and read heads that are suspended by a suspension arm above the rotating disk and an actuator that swings the suspension arm to place the read and write heads over selected circular tracks on the rotating disk. The read and write heads are directly mounted on a slider that has an air bearing surface (ABS). The suspension arm biases the slider into contact with the surface of the disk when the disk is not rotating but, when the disk rotates, air is swirled by the rotating disk adjacent the ABS to cause the slider to ride on an air bearing a slight distance from the surface of the rotating disk. When the slider rides on the air bearing the write and read heads are employed for writing magnetic impressions to and reading magnetic impressions from the rotating disk. The read and write heads are connected to processing circuitry that operates according to a computer program to implement the writing and reading functions.\nThe write head includes a coil layer embedded in first, second and third insulation layers (insulation stack), the insulation stack being sandwiched between first and second pole piece layers. A gap is formed between the first and second pole piece layers by a gap layer at an air bearing surface (ABS) of the write head. The pole piece layers are connected at a back gap. Current conducted to the coil layer induces a magnetic field across the gap between the pole pieces. This field fringes across the gap at the ABS for the purpose of writing the aforementioned magnetic impression in tracks on moving media, such as in circular tracks on the aforementioned rotating disk.\nIn recent read heads a spin valve sensor is employed for sensing magnetic fields from the rotating magnetic disk. The sensor includes a nonmagnetic conductive layer, hereinafter referred to as a spacer layer, sandwiched between first and second ferromagnetic layers, hereinafter referred to as a pinned layer, and a free layer. First and second leads are connected to the spin valve sensor for conducting a sense current therethrough. The magnetization of the pinned layer is pinned perpendicular to an air bearing surface (ABS) of the head and the magnetic moment of the free layer is located parallel to the ABS but free to rotate in response to external magnetic fields. The magnetization of the pinned layer is typically pinned by exchange coupling with an antiferromagnetic layer.\nThe thickness of the spacer layer is chosen to be less than the mean free path of conduction electrons through the sensor. With this arrangement, a portion of the conduction electrons is scattered by the interfaces of the spacer layer with the pinned and free layers. When the magnetizations of the pinned and free layers are parallel with respect to one another, scattering is minimal and when the magnetizations of the pinned and free layers are antiparallel, scattering is maximized. Changes in scattering alter the resistance of the spin valve sensor in proportion to cos .theta., where .theta. is the angle between the magnetizations of the pinned and free layers. In a read mode the resistance of the spin valve sensor changes proportionally to the magnitudes of the magnetic fields from the rotating disk. When a sense current is conducted through the spin valve sensor, resistance changes cause potential changes that are detected and processed as playback signals.\nA spin valve sensor is characterized by a magnetoresistive (MR) coefficient that is substantially higher than the MR coefficient of an anisotropic magnetoresistive (AMR) sensor. For this reason a spin valve sensor is sometimes referred to as a giant magnetoresistive (GMR) sensor. When a spin valve sensor employs a single pinned layer it is referred to as a simple spin valve. A spin valve is also know as a top or bottom spin valve depending upon whether the pinning layer is at the top (formed after the free layer) or at the bottom (before the free layer). A pinning layer in a bottom spin valve is typically made of nickel oxide (NiO).\nAnother type of spin valve sensor is an antiparallel (AP) spin valve sensor. The AP pinned spin valve sensor differs from the simple spin valve sensor, described above, in that the pinned layer of the AP pinned spin valve sensor comprises multiple thin layers, which are collectively referred to as an antiparallel (AP) pinned layer. The AP pinned layer has a ruthenium (Ru) spacer layer sandwiched between first and second ferromagnetic thin layers. The first ferromagnetic thin layer has its magnetic moment oriented in a first direction by exchange coupling to the antiferromagnetic pinning layer. The second ferromagnetic thin layer is immediately adjacent to the free layer and is antiparallel coupled to the first thin layer because of the minimal thickness (in the order of 8 .ANG.) of the spacer layer between the first and second ferromagnetic thin layers. The magnetic moment of the second ferromagnetic thin layer is oriented in a second direction that is antiparallel to the direction of the magnetic moment of the first ferromagnetic layer.\nThe AP pinned layer is preferred over the single layer pinned layer. The magnetic moments of the first and second layers of the AP pinned layer subtractively combine to provide a net pinning moment of the AP pinned layer. The direction of the net moment is determined by the thicker of the first and second thin layers. The thicknesses of the first and second thin layers are chosen to reduce the net moment. A reduced net moment equates to a reduced demagnetization (demag) field from the AP pinned layer. Since the antiferromagnetic exchange coupling is inversely proportional to the net pinning moment, this increases exchange coupling between the first ferromagnetic film of the AP pinned layer and the pinning layer. The high exchange coupling promotes higher stability of the head. When the head encounters elevated thermal conditions caused by electrostatic discharge (ESD) from an object or person, or by contacting an asperity on a magnetic disk, the blocking temperature (temperature at which magnetic spins of the layer can be easily moved by an applied magnetic field) of the antiferromagnetic layer can be exceeded, resulting in disorientation of its magnetic spins. The magnetic moment of the pinned layer is then no longer pinned in the desired direction. A reduced demag field also reduces the demag field imposed on the free layer which promotes a symmetry of the read signal. The AP pinned spin valve sensor is described in commonly assigned U.S. Pat. No. 5,465,185 to Heim and Parkin which is incorporated by reference herein.\nThe first and second ferromagnetic layers of the AP pinned spin valve sensor are typically made of cobalt (Co). Unfortunately, cobalt has high coercivity, high magnetostriction and low resistance. When the first and second ferromagnetic layers are formed they are sputtered deposited in the presence of a magnetic field that is oriented perpendicular to the ABS which sets the easy axis (e.a.) of the ferromagnetic films perpendicular to the ABS. During operation of the head the AP pinned layer is subjected to extraneous magnetic fields that have components parallel to the ABS, such as components of the write field. These extraneous fields, combined with heating of the pinning layer, can cause the pinning layer to lose its pinning strength (exchange coupling) and allow the magnetic moments of the ferromagnetic layers to switch from being perpendicular to the ABS to some other direction. If the coercivity of the ferromagnetic films is higher than the exchange field that urges the magnetic moments of the ferromagnetic layers back to their original positions the magnetic moments of the ferromagnetic layers will remain in the wrong direction. This renders the read head inoperable.\nCobalt (Co) has a high negative magnetostriction. The negative sign determines the direction of any stress induced anisotropy. When a magnetic head is lapped, which is a grinding process, nonuniform compressive stresses occur in the layers of the sensor. Because of the magnetostriction and the stresses the cobalt (Co) ferromagnetic films acquire a stress induced anisotropy that is parallel to the ABS. This is the wrong direction. The stress induced anisotropy may rotate the magnetic moment of the first and second ferromagnetic layers of the AP pinned layer to some extent from perpendicular to the ABS in spite of the exchange coupling field tending to maintain the perpendicular position. This condition can cause read signal asymmetry.\nThe low resistance of the cobalt (Co) ferromagnetic films of the AP pinned layer causes a portion of the sense current to be shunted past the free and spacer layers. This causes a loss of read signal.\nEfforts continue to increase the spin valve effect of GMR heads. An increase in the spin valve effect equates to higher bit density (bits/square inch of the rotating magnetic disk) read by the read head. Promoting read signal symmetry is also a consideration. This is accomplished by reducing the magnetic influences on the free layer. A search still continues to lower the coercivity, substantially eliminate magnetostriction and increase the resistance of some of the critical layers of the spin valve sensor."}
{"text": "The present invention relates to a device and to a method for the control of the air-conditioning system of a vehicle, in particular a motor vehicle such as an automobile, a bus, etc.\nAs is known, the air-conditioning systems of motor vehicles are typically provided with a closed-loop cooling circuit equipped with an evaporator and a control system capable of regulating the temperature of the air introduced into the passenger compartment.\nFIG. 1 shows a control system 1 of a cooling circuit 2, which is traversed by the coolant and comprises in succession: an evaporator 3; a compressor 4, which is designed to take in, at a certain intake pressure, the coolant in the vapour phase from the evaporator 3 so as to obtain a control of the temperature of the air downstream of the evaporator 3 itself; a condenser 5, designed to receive the coolant in the vapour phase from the compressor 4; and an expansion valve 6 designed to receive the coolant in the liquid phase from the condenser 5 to supply it in dual-phase (i.e., vapour phase and liquid phase) to the evaporator 4 itself.\nIn particular, the compressor 4 is constituted by a compressor with externally controlled variable displacement, on which it is possible to operate by causing the displacement (defined as the working volume, where the coolant is compressed) to vary as the thermal load acting on the air-conditioning system varies.\nThe possibility of varying the displacement of the compressor is obtained via the electromagnetic regulation valves (not illustrated), which are driven via an external control signal SC and are designed to control the device that enables modulation of the displacement of the compressor 4.\nThe control system 1 moreover comprises a control device 7, which is able to generate the control signal SC of the compressor 4 in such a way as to control the temperature of the air downstream of the evaporator 3 as a function of the deviation between a reference temperature TREF set by the user by means of an externally controlled selector device 9 and an effective temperature TMIS indicating the temperature of the air present downstream of the evaporator 3. The effective temperature TMIS can be measured using a temperature sensor 10, set downstream of the evaporator 3.\nThe control device 7 comprises an adder block 11 having a first input designed to receive the reference temperature TREF from the selector device 9, a second input designed to receive the effective temperature TMIS from the temperature sensor 10, and an output supplying a temperature error er, given by the difference between the reference temperature TREF and the effective temperature TMIS.\nThe control device 7 moreover comprises a control block 12, which is designed to receive at input the temperature error er and a set of measurement parameters, such as for example Te (external temperature), RPM (engine r.p.m.), and RH (relative humidity) correlated to the exogenous disturbance, and supplies at output, according to the latter, the control signal SC.\nIn detail, the control block 12 comprises a compensating network of a proportional-integral (PI) type (not illustrated) and generates a control signal SC corresponding to a pulse-width modulation (PWM) signal, which drives the electromagnetic valves for regulating the stroke of the pistons, thus determining control of the displacement of the compressor 4. In the case in point, the regulation of the displacement determines a control of the intake pressure of the compressor 4 and, consequently, an indirect control of the temperature of the air downstream of the evaporator 3.\nIt is moreover known that, in the air-conditioning systems described above, the expansion valve 6 for supplying the coolant to the evaporator 3 is an internally controlled device, operation of which is completely independent of the control implemented on the compressor 4 by the control device 7.\nThe complete independence existing between the two controls determines, in certain limit conditions of operation of the compressor 4, a discordance in the control of some parameters that characterize operation of the cooling circuit 2, such as for example the intake pressure of the compressor 4 and the temperature of the air downstream of the evaporator 3, in this way causing a condition of instability of the air-conditioning system. In the case in point, during its operation, the expansion valve 6 generates a temperature disturbance ΔTEVAP of an oscillatory type, which alters the temperature of the air downstream of the evaporator 3, and leads, in certain conditions, to instability of the control.\nIn fact, said temperature disturbance ΔTEVAP determines an increase of the intake pressure of the compressor 4, which in certain limit conditions exceeds a threshold delimiting the condition of stability of the air-conditioning system, consequently causing a series of oscillations of the flow rate of the coolant, and of the temperature of the air downstream of the evaporator 3. The generation of said oscillations, generally referred to with the term “hunting phenomenon”, represents a major drawback in the air-conditioning systems described above in so far as it has a negative effect both on the capacity of minimizing the consumption of the air-conditioning system, and on the thermal comfort of the passenger compartment of the vehicle."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical coupling device; and, more particularly, to an optical coupling device having a silicon optical bench (SiOB) and an optical fiber with an angled end face.\n2. Description of Related Arts\nAs a high-speed Internet network is widely distributed and subscribers to this high-speed Internet network service rapidly increase, the amounts of data traffic also increase. Thus, a telecommunication system that uses electric signals confronts a limitation to meet the above demand. For this reason, an optical communication system is started to be the alternative of the conventional telecommunication system since it is able to transmit large amounts of data. Optical fibers, optical transmitter modules and optical receiver modules are representative optical components used in the optical communication system. For the optical transmitter and receiver modules high-speed, small size, low cost, and an ease of mass-production are key factors in the view of commercialization.\nSince a typical optical module using a TO-can packaged optical device has problems in a high-speed operation and a large volume, an optical module using a silicon optical bench (hereinafter referred as to SiOB) technology has been currently developed and commercialized. Herein, the SiOB provide a function for an optical coupling of an optical fiber to a laser by including solder-bumps and a V-shaped groove. A laser chip is aligned and mounted on the SiOB. The V-shaped groove is the place in which an optical fiber is aligned and fixed, and the solder bumps are used for alignment and bonding of a laser chip.\nOptical modules using the SiOB have advantages of reduced assembly costs because the optical fiber and the optical active elements such as a laser diode and a photodiode can be aligned by a passive alignment technique such as a conventional flip-chip bonding technique. Also, the optical modules can be minimized due to the fact that plural of optical fiber and arrayed optical active elements can be integrated onto one substrate. FIG. 1 is a diagram showing an optical coupling device using an edge emitting laser and a SiOB structure in accordance with a prior art. (Refer to references G. C. Joo, S. H. Lee, K. S. Park, N. Hwang, J. S. Choi and M. K. Song, “Bidirectional optical coupling of transceiver chip for subscribers”, Electron. Lett., vol. 34, no. 24, November 1998 and M. H. Choi, H. J. Koh, E. S. Yoon, K. C. Shin and K. C. Song, “Self-Aligning Silicon Groove Technology Platform for the Low Cost Optical Module”, IEEE ECTC, pp. 1140–1144, 1999.)\nWith reference to FIG. 1, a groove 11 is formed on a surface of a SiOB 10. An optical fiber 16 is located on the groove 11. Also, on top of the SiOB 10, a core 17 of the optical fiber 16 is located at a height identical to a height at which an active area of an edge emitting laser 14 is formed. The edge emitting laser 14 is mounted on the SiOB 10 by a metal pad 12 and a solder 13. Typically, a monitoring photodiode 15 is allocated in an opposite to the optical fiber 16.\nA light emitted from the edge emitting laser 14 is transmitted 18 through the core 17 of the optical fiber 16 and the monitoring photodiode 15 receives the rear emitted light 19 of the laser 14.\nFIG. 2 is a diagram showing an optical coupling device with the use of a surface emitting laser and an optical fiber having an angled end face in accordance with a prior art (Refer to the U.S. Pat. No. 6,389,202 and the U.S. Pat. No. 6,315,464).\nWith reference to FIG. 2, one end face of an optical fiber 30 to which light is incident is formed as to have an angle of about 45° with respect to an optical fiber axis, and a mirror plane 32 coated with metal such as Au or Al is formed on the angled end face of the optical fiber 30. A vertical cavity surface emitting laser (VCSEL) 33 is allocated at a bottom portion of the angled end face of the optical fiber 30. In this case, a laser light emitted from the VCSEL 33, that is, an incident light IL incident to the optical fiber 30 is reflected at the mirror plane 32 and transmitted through the core 31 of the optical fiber 30. Herein, the light transmitted through the core 31 of the optical fiber 30 is called TL.\nHowever, in the case of applying this optical coupling structure, a reflection light RL is inevitably generated as a portion of the incident light IL emitted from the VCSEL is reflected at end facean incident plane of an optical fiber 30. The reflected light is incident to the VCSEL again. As a result, the VCSEL 33 is degraded due to the reflection light RL."}
{"text": "The present invention relates generally to improved circuits and methods for generating clock signals by multiplying a fundamental clock frequency, and more particularly to substantially increasing the accuracy of the multiplied clock frequency above the accuracy achievable in the prior art.\n“Prior Art” FIG. 1 shows a clock generation circuit 1 which includes a typical relaxation oscillator 2 that is believed to be representative of the closest prior art. Relatively “tight” frequency specifications are required for some integrated circuits including internal clock generators based on internal relaxation oscillators. For example, relatively tight frequency specifications are required for internal clock signals in integrated circuits which are designed to communicate over an RS232 interface. In the manufacture of such an integrated circuit, it may be desirable to perform a laser trimming operation on certain resistors in the integrated circuit at room temperature to produce a desired relaxation oscillator frequency. However, after the laser trimming operation the relaxation oscillator frequency nevertheless may vary significantly over a specified temperature range (e.g., minus 50 degrees Centigrade to plus 125 degrees Centigrade). There may be a need to ensure that the laser trimming operation at room temperature results in the relaxation oscillator frequency being within the specified range over the entire allowable temperature range. This has been difficult to achieve for internal clock generation circuits based on prior art relaxation oscillator circuits.\nReferring to FIG. 1, clock generator circuit 1 includes a relaxation oscillator 2 including a current source 11 which supplies a current I1 to the sources of P-channel transistors M1 and M2, the gates of which are controlled by a logic circuit 10. The drain of transistor M1 is connected by conductor 3 to capacitor C1, switch S2, and a comparator 13. The drain of transistor M2 is connected to capacitor C2, switch S1, and comparator 12. Comparators 12 and 13 receive a reference voltage Vref that is produced by a current source Iref flowing through a resistor having a resistance R1.\nThe outputs of comparators 12 and 13 are connected to inputs of an OR gate 14, the output of which clocks a flip-flop 15. The output 16 of flip-flop 15 is connected to the inputs of a inverter 18 and a non-inverting buffer 20 and to an input of logic circuit 10. Logic circuit 10 produces clock signals φ1, φ2, φ1, and φ2 in response to the signal produced by flip-flop 15, where φ2 can be equal to φ1 and φ1 can be equal to φ2.\nIf transistor M1 is on, transistor M2 is off, switch S1 is closed, and switch S2 is open, causing the current I1 to flow through transistor M1 and gradually charge up capacitor C1, producing the ramp section “A1” of V3 as shown during which the clock signal CLK in the timing diagram of FIG. 2 is at the level P1. When V3 exceeds reference voltage Vref, the output 36 of comparator 12 switches from a “0” level to a “1” level after a propagation delay Tprop, causing a “1” level to be produced at the output of OR gate 14. This causes flip-flop 15 to change state, causing buffer 20 to produce the low level N1 of clock signal CLK in FIG. 2 and causing logic circuit 10 to switch φ1 to a high level and to switch φ2 to a low level\nThat turns transistor M1 off, turns transistor M2 on, closes switch S2, and opens switch S1. The closing of switch S2 discharges capacitor C1, producing transition B1 of voltage signal V3. The current I1 flows through transistor M2 and gradually charges up capacitor C2, producing the ramp section “A2” of signal V4. The operation continues similarly to that described above, and when V4 exceeds reference voltage Vref, the output 37 of comparator 13 switches from a “0” level to a “1” level, and a “1” level then produced at the output of OR gate 14 causes flip-flop 15 to change state, reversing the levels of φ1 and φ2. This discharges capacitor C2 to ground, causing transition “B2” of voltage signal V4 and also causing the transition of clock signal CLK to the level P2. As the foregoing operation is repeated, the level of CLK changes and the levels of φ1 and φ2 change each time the ramp portions of V3 and V4 exceed Vref.\nThe period of oscillation TOSC of relaxation oscillator 2 is given by the expressionTOSC=(2*C*Vref)/I1+Tprop.   Eq. (1)\nThe value of Tprop has a great influence on the ability of relaxation oscillator 2 to operate at very high clock frequencies (e.g., several hundred megahertz) over a typical expected temperature range (e.g., −50 degrees Centigrade to +125 degrees Centigrade) because Tprop can vary significantly over that range.\nFor clock signal periods requiring an accuracy of less than, for example 2%, over a predetermined typical temperature range and power supply range, the variance of the propagation delay Tprop can limit the maximum clock frequency obtainable. For example, a deviation of 1 nanosecond in Tprop can produce a corresponding 1% deviation in a 10 MHz nominal frequency of relaxation oscillator 2 and hence in the frequency of CLK. However, since the propagation delay Tprop is not totally dependent on the frequency of relaxation oscillator 2, if in this example relaxation oscillator 2 oscillates at a frequency higher than 10 MHz, any deviation of the propagation delay Tprop over the expected temperature range results in a correspondingly greater percentage deviation than 1% in the corresponding period of relaxation oscillator 2.\nThus, there is an unmet need for a clock generator circuit based on a relaxation oscillator that provides increased accuracy of the generated clock signal.\nThere also is an unmet need for a circuit and technique for multiplying the frequency of a signal produced by a relaxation oscillator without having to increase the current supplied by a current source in the relaxation oscillator or varying the value of capacitors such as capacitors C1 and C2 in FIG. 1.\nThere also is an unmet need for a circuit and technique for multiplying the frequency of a signal based on a clock signal derived from a relaxation oscillator without having to increase the frequency of the relaxation oscillator.\nThere also is an unmet need for an analog technique and circuit for multiplying an output signal produced by a relaxation oscillator to obtain a signal having a frequency which is a multiple of the frequency of the output signal produced by the relaxation oscillator without reducing overall circuit performance of the relaxation oscillator and the analog multiplying circuit over an expected temperature range.\nThere also is an unmet need for an analog technique and circuit for multiplying an output signal produced by a relaxation oscillator to obtain a signal having a frequency which is a multiple of the frequency of the output signal produced by the relaxation oscillator without unacceptably limiting the maximum achievable frequency of the multiplied output signal.\nThere also is an unmet need for an analog technique for obtaining signals within certain phase specifications or time delay specifications of a fundamental relaxation oscillator output."}
{"text": "This invention relates to functionalized cubic liquid crystalline phases and methods for their preparation and use. More specifically, this invention relates to functionalized cubic liquid crystalline phase materials that have properties tailored to specific uses.\nMuch of the interest in bicontinuous cubic phase liquid crystals is a consequence of their unique structure. They are composed of mixtures of lipid and water arranged into bilayers. The bilayers, in turn, are twisted into a periodic, three-dimension structure that minimizes the energy associated with bending the bilayers (i.e. minimize curvature energy). See Hyde, S., Andersson, S., Larrson, K., Blum, Z., Landh, T., Lidin, S., Ninham, B. W., The Language of Shape, Elsevier Press, New York, 1997. These structures are xe2x80x98honeycombedxe2x80x99 with bicontinuous domains of water and lipid that is reminiscent of an organic zeolite or highly-structured microemulsion. As such the structure can simultaneously accommodate water-soluble, lipid-soluble and amphiphilic molecules and provide pathways for diffusion of water-soluble and lipid-soluble materials. While there have been a number of proposed cubic phases, there are three recognized bicontinuous liquid crystals structures: Pn3m (D-surface), Ia3d (G-surface), and Im3m (P-surface). See Luzzati, V., Vargas, R., Mariani, P., Gulik, A., Delacroix, H., J. Mol. Biol., 1993, 229, 540-551. These structures can be difficult to express in rigorous mathematical terms. However, if expressed in terms of nodal surfaces, structure and shape can be approximated. See von Schnering, H. G., Nesper, R. Z., Phys. B-Condensed Matter, 1991, 83, 407-412. The phase behavior of a broad range of monoglycerides has been documented and modifications to the phase behavior have been defined. See Qiu, H., Caffrey, M., xe2x80x9cThe phase diagram of the monoolein/water system: metastability and equilibrium aspectsxe2x80x9d, Biomaterials, 1999, 21(3), 223-234. Accordingly, monoolein-based bicontinuous cubic liquid crystal phase have good temperature stability, high internal surface area, gel-like viscosity, relative insensitivity to salt and solvent compositions, and use low cost raw materials which make them practical for commercial applications. Monoolein naturally exhibits Pn3m and Ia3d, with Im3m present with the addition of proteins. See Rummel, G., Hardmeyer, A., Widmer, C., Chiu, M. L., Nollert, P., Locher, K. P., Pedruzzi, I., Landau, E. M., Rosenbusch, J. P., J. Structural Biology, 1998, 121, 82-91.\nCubic phase liquid crystals have been used in gel, dispersion and precursor form. xe2x80x98Gelsxe2x80x99 are mixtures that contain a majority of the cubic phase liquid crystal. It is common for either mixture to exclusively contain cubic liquid crystal phase. Applications for these gels can range from drug delivery vehicles (See Shah, J. C., Sadhale, Y., Chilukuri, D. M., Adv. Drug Delivery Rev., 2001, 47(2-3), 229-250), to a matrix in which membrane proteins can be crystallized (See Landau, E., Rosenbusch, J., Proc. Natl. Acad. Sci. USA., 1996, 93(25), 14532-14535), or in which mesoporous nanoparticles can be formed (See Cruise, N., Jansson, K., Holmberg, K., J. Colloid Interface Sci., 2001, 241(2), 527-529).\nNielsen, WO 98/47487, discloses compositions of bio-adhesive liquid crystal gels, including the cubic phase liquid crystals and precursors. Compostions include an active, a cubic phase forming lipid, and a structurant that is added without changing the structure of the liquid crystal. These compositions do not disclose the use of hydrotropes to form liquid crystals.\nEngstrom et al., U.S. Pat. No. 5,753,259, discloses a composition and method of use of liquid crystal gels, including cubic phase liquid crystals, for controlled release applications. The disclosed gels are fabricated from a mixture of lipid, solvent, and bioactive materials including nucleic acids. However, these gel compositions do not utilize hydrotropes.\nxe2x80x98Dispersionsxe2x80x99 are particles of cubic liquid crystalline phase material that are often submicron in size. Particles are generally dispersed in a liquid medium and are often termed Cubosomes. Cavitating a mixture of lipid and liquid generally makes dispersions of cubic phase liquid crystals. This requires high pressures and numerous passes before homogeneous nanoparticle dispersions are produced (See Ljusberg-Wahren, H., Nyberg, L., Larsson, K., Chimica Oggi, 1996, 14, 40-43). Cubosomes have distinct practical advantages over vesicles and liposomes because cubosomes are an equilibrium phase (See Laughlin, R. G. Colloids and Surfaces A, 1997, 128, 27-38). Cubosomes also possess much greater internal surface area than vesicles or liposomes and are more resilient against degradation.\nAnderson, WO 99/12640, and Landh et al., U.S. Pat. No. 5,531,925, disclose cubic phase compositions and preparations for delivery and uptake of active agents. The particles comprise a center containing liquid crystalline or liquid material and an exterior of solid particulate. The composition of the liquids comprise a lipid and polar solvent without the sue of hydrotropes.\nxe2x80x98Precursorsxe2x80x99 are mixtures that are not cubic phase liquid crystals but form cubic phase liquid crystals as a consequence of a stimulus. Precursors can be used to dispense a mixture in a form that readily flows, but spontaneously converts to a more viscous liquid crystal with the stimulus at a target location. This is applicable to treatments for periodontal disease (See Norling, Tomas, Lading, Pia, Engstroem, Sven, Larsson, Kare, Krog, Niels, Nissen, Soeren Soe, J. Clin. Periodontol, 1992, 19(9, Pt. 2), 687-92. Larson et al., U.S. Pat. No. 5,196,201, discloses the preparation and composition of precursors used as implants to treat aliments such as the repair of bone tissue. These precursors are composed of a water-based liquid, lipid, and optionally a triglyceride mixed to form a more concentrated L2 or D phase, which flows more readily, and converts to cubic phase upon the addition of water. Leng et al., U.S. Pat. No. 5,593,663, discloses combinations and preparations of antiperspirant, which uptake sweat upon application to form a viscous liquid crystalline phase, including cubic phase. However, neither of these materials contains functionalization materials.\nCubic liquid crystalline phase materials are limited in use due to restriction of their natural, or unmodified, properties. For example, the natural properties of cubic phases limit the ability to solubilize active ingredients. In fact, broad classes of actives do not effectively load (or subsequently release) because the cubic phase lacks specific interaction with the loaded active. If the active is modified to effectively load in the cubic phase, it may lose its effectiveness. Further, there are no commercially convenient ways to provide specific targeting or enhanced deposition of actives from cubic phase. Finally, there are no cubic phases suitable for xe2x80x98on demandxe2x80x99 applications. xe2x80x9cOn demandxe2x80x99 refers to changes in the properties of cubic phase as a consequence of some stimulus, such as change in pH. As a result, a technique is needed to modify the cubic phase and significantly increase the utility of cubic phase.\nA cubic liquid crystalline phase precursor comprising (A) a hydrotrope, (B) an amphiphile capable of forming a cubic liquid crystalline phase, (C) an optional solvent, and (D) an additive selected from the group consisting of an anchor, a tether, and combinations thereof, wherein ingredients (A), (B), (C), and (D) are present in mass fractions relative to each other such that 1.0=a+b+c+d, wherein a is the mass fraction of ingredient (A), b is the mass fraction of ingredient (B), c is the mass fraction of ingredient (C), and d is the mass fraction of ingredient (D), and wherein 1.0 greater than a greater than 0, 1.0 greater than b greater than 0, 1.0 greater than cxe2x89xa70, 1.0 greater than d greater than 0; and with the proviso that a, b, c, and d do not fall within a cubic liquid crystalline phase region on a phase diagram representing phase behavior of ingredients (A), (B), (C), and (D).\nA bulk cubic liquid crystalline gel comprising (A) a hydrotrope, (B) an amphiphile capable of forming a cubic liquid crystalline phase, (C) a solvent, and (D) an additive selected from the group consisting of an anchor, a tether, and combinations thereof, wherein ingredients (A), (B), (C), and (D) are present in mass fractions relative to each other such that 1.0=a+b+c+d, wherein a is the mass fraction of ingredient (A), b is the mass fraction of ingredient (B), c is the mass fraction of ingredient (C), and d is the mass fraction of ingredient (D), and wherein 1.0 greater than a greater than 0, 1.0 greater than b greater than 0, 1.0 greater than c greater than 0, 1.0 greater than d greater than 0; and with the proviso that a, b, c, and d fall within a cubic liquid crystalline phase region on a phase diagram representing phase behavior of ingredients (A), (B), (C), and (D).\nA disperison of cubic liquid crystalline gel particles comprising (A) a hydrotrope, (B) an amphiphile capable of forming a cubic liquid crystalline phase, (C) a solvent, and (D) an additive selected from the group consisting of an anchor, a tether, and combinations thereof, wherein ingredients (A), (B), (C), and (D) are present in mass fractions relative to each other such that 1.0=a+b+c+d, wherein a is the mass fraction of ingredient (A), b is the mass fraction of ingredient (B), c is the mass fraction of ingredient (C), and d is the mass fraction of ingredient (D), and wherein 1.0 greater than a greater than 0, 1.0 greater than b greater than 0, 1.0 greater than c greater than 0, 1.0 greater than d greater than 0; and with the proviso that a, b, c, and d fall within a region representing cubic liquid crystalline phase in combination with at least one other phase on a phase diagram representing phase behavior of ingredients (A), (B), (C), and (D), with the proviso that the dispersion has a form of functionalized cubic liquid crystalline gel particles dispersed in the other phase."}
{"text": "This invention relates to cooking utensils, namely oven basters, and more particularly to an oven baster and a cleaning brush therefor.\nA cooking utensil that is useful in the kitchen is the oven baster. Such a utensil is used to apply liquids to foods, typically meat or poultry, that are being baked or roasted in an oven. The liquid is applied periodically during cooking to prevent excessive drying of the food or to add flavorings. Typically, prior oven basters have included a tube of relatively large diameter having a tip portion of reduced diameter at one end and a compressible, self-expanding bulb at the other end. With the bulb compressed, the tip portion is placed in a reservoir of basting fluid. Pressure on the bulb is then released, allowing the bulb to self-expand and draw fluid into the tube. By positioning the tip portion over the food to be basted and again compressing the bulb, the fluid is expelled from the tube through the tip portion and onto the food.\nThe liquids that are typically dispensed using an oven baster often include solid particles or other substances that adhere to the interior surface of the basting tube, or that are large enough to become lodged in the tip portion of the tube. It is necessary to clean the interior of the tube after use. This can be accomplished by disassembling the tube and bulb and using a bottle-type brush to clean the interior of the basting tube. Such brushes, if sized for the relatively large diameter of the reservoir portion of the tube, are too large to fit within the smaller diameter tip portion. Either the tip portion must be cleaned with water and determent alone, or some other utensil of small diameter must be used to clean or dislodge matter from the tip portion, which is inconvenient at best.\nIt would be desirable to provide a cleaning brush for an oven baster that alleviates the difficulty of cleaning the small diameter tip portion portion of the basting tube.\nPrior oven basters have employed tubes constructed of various materials, such as metal, glass and opaque or translucent plastic. Glass and metal are resistant to heat deformation. Glass also has the advantage of being transparent, which favors inspection and measurement of basting fluid within the tube, and also verifying cleanliness of the interior of the tube. Glass is quite easily broken during storage or use, however. Opaque or translucent plastic is resistant to breakage, but is not as heat resistant as glass and does not have the advantages of transparency.\nIt would be desirable to provide an oven baster having a tube that is transparent, resistant to breakage, and resistant to heat.\nThese and other advantages are provided by the present invention.\nThe present invention, according to one aspect, includes a baster having a tube with a tip portion and a reservoir portion. The tube is comprised of a substantially transparent polymeric material. A compressible bulb is in communication with the reservoir portion of the tube. The bulb is resiliently self-expandable to induce reduced air pressure in the tube relative to ambient air pressure.\nAccording to another aspect of the present invention, a cleaning brush for a tube includes a handle having a proximate end and a distal end. A longitudinal shank extends distally from the distal end of the handle, and plurality of bristles extend substantially transversely from the shank. An elongate cleaning probe extends distally from the distal end of the handle.\nAccording to yet another aspect of the present invention, a baster and cleaning brush set is provided. Included is a baster having a tube with a reservoir portion and a tip portion, and a compressible bulb in communication with the reservoir portion of the tube. The bulb is resiliently self-expandable to induce reduced air pressure in the tube relative to ambient air pressure. Also included is a cleaning brush for the tube having a handle with a proximate end and a distal end. A longitudinal shank extends distally from the distal end of the handle. A plurality of bristles extend substantially transversely from the shank. An elongate cleaning probe extends distally from the distal end of the handle.\nOther features and advantages of the present invention will be apparent from the following descriptions of a preferred embodiment, made with reference to the drawings."}
{"text": "There are multiple technologies that can be used in the access network but as it is known, connection paths comprising more optical fiber (relative to copper) will typically yield greater data rates, so a network comprising fiber all the way to customer premises (FTTP) which may take the form of a passive optical network (PON) is the most preferred solution in terms of transmission capacities and speeds.\nAs demand for capacity of the fiber access networks increases, eventually multiple wavelengths are used between the access network head-end (OLT) and a remote terminal at some location near the customer premises depending on the type of FTTx system being used (e.g. FTTP, FTTC fiber to the cabinet, FTTDP fiber to the distribution point).\nA schematic depiction of a conventional PON (10) is depicted in FIG. 1, in which an OLT (14) includes one or more core network linecards each of which has a plurality of core network optical ports, and one or more linecards (40) each of which has a plurality of optical PON line terminals or ports (typically 4, 8, or 16, also not shown). Data is transmitted to the customer end via a PON port or line terminal (LT) along a connection established over an optical distribution network (ODN) in the access network. This comprises a primary PON feeder fiber connection (34), an optical splitter (42) and a number of optical distribution fiber links (46) to an optical network unit (ONU) (also known as an optical network terminal ONT), which can be located at different points in the network close by or inside the customer premises (32). For ease of description, FIG. 1 depicts only one PON and the relevant section of the network, but it is known that an OLT can support many PON systems via multiple ODNs and serve many more ONUs.\nNetworks are susceptible to equipment or plant damage or failure and have limited resilience in the event of equipment or plant failure. The adverse effects can be made worse by the network branching architecture. For example, if the optical feeder fiber link (34) connecting the OLT to an optical splitter (42) is broken, or if a PON linecard serving multiple ONUs was to develop a fault, then all of the customers served by that PON port or the entire PON linecard will lose service. Thus, a problem with the OLT (14) could lead to a total loss of service, or at the least reduction in service quality, for all of the customers served by the affected OLT.\n“Dual parenting” (also known as “dual homing”) is a network design approach addressing the need to provide operational protection or resilience in optical networks in the event of loss of service to customers. Here, a customer is connected to head ends in two different exchange buildings to increase the resilience of the network.\nWO2009/050459 describes a dual parenting implementation in which a PON is configured with two virtual local area network (VLAN) connections between a switch and the customer terminal, using two different communication paths provided by dual parenting. In operation, one of the VLAN connections is used to transmit traffic under “normal” network conditions and the other provides a back-up path in the event of problems (such as cable breaks, equipment malfunction, or damage or destruction of the exchange building housing the primary OLT) affecting operation of the first connection. For the purpose of this description, “normal” network conditions refers to the situation where the network or part thereof is operating as expected in the absence of damage, failure or faults affecting network plant and equipment affecting the network or section under consideration. The provision of redundancy in this way at the network layer and the service layer enables a faster customer reconnection in the event of network plant or equipment failure. In this document, this approach is also referred as “protecting” the network or providing it with operational resilience.\nAnother dual parented implementation is described in EP13250023.2, in which a dual parenting network topology is deployed in a fiber to the cabinet (FTTC) network, in which the cabinets are connected to the exchange via point-to-point optical fiber links. Here, the standby network has a point-to-multipoint layout.\nProviding redundancy in this manner significantly increases the number of standby line cards, network plant and equipment needed and thus the overall cost of building the network.\nRedundant plant and equipment are typically not used to their fullest extent, if at all, under normal network operation conditions until there is a need to switch operations away from the primary OLT. Failure events are rare occurrences and so significant capital and operation expense is incurred in providing a standby OLT and its associated PON ports.\nOther protection schemes are known in the prior art, such as that described in US2012/0251108, which discloses an optical network comprising two multiplexer/demultiplexers (mux/demux), the first being comprised in an OLT, and the second forming a wavelength routed wavelength division multiplexed PON. The network further comprises working feeder fibers and protection backup feeder fibers connected to the mux/demux units comprising an N×N or 1×N arrayed wavelength gratings (AWGs). The output ports of the second mux/demux are allocated as a working path output port or a protection path output port depending on: an identity of a second mux/demux, an identity of the output port, an identity of a free spectral range (FSR) of the first mux/demux, the number of working fibers, the number of input ports of the first mux/demux, and an identity of the backup feeder fiber.\nThe problem is also addressed in “Flexible and scalable PON protection architecture using N:M redundancy toward next generation access network” (Mitsui T; NTT Access Network Service Syst. Labs, NTT Corp, Yokosuka Japan; Sakamoto T; Hara K; Yoshimoto N, The 17th Asia-Pacific Conference on Communication, 2011), which describes a system having an optical switch between the splitters situated between the ONUs and the OLTs. If one of the OLTs fails, the fibers directed to the failed OLT are switched to another OLT, using an optical switch in order to implement an N:M protection (i.e. N PONs protected by M OLT ports). This application is restricted to use in FTTC networks only, where there is already an active element in the field, i.e. the DSLAM. The idea is not usable in the FTTP technology where the fiber plant and equipment is passive. It is also not suited for use in a dual parenting resilience scheme where the working and protecting OLT's are sited in remote locations.\nYet further approaches seek to protect the fibers between the remote node and the ONUs in specific configurations, such as those described in US2006104638 and MY140726."}
{"text": "Field of the Invention\nThe invention relates to a modular printing machine system for printing on sheets, including a first printing machine of satellite construction type having a central first impression cylinder and at least four printing devices assigned thereto, a second printing machine having a second impression cylinder, and a coupling device for coupling the printing machines to one another for in-line operation thereof.\nThe published German Patent Document DE 197 43 770 A1, which corresponds to U.S. Pat. No. 6,101,944, describes such a printing machine system, which includes a satellite printing-unit group as the first printing machine, and imprinting-unit group as the second printing machine and a connecting element as a coupling device. A point of separation in the coupling runs between a transfer drum belonging to the satellite printing-unit group and a transfer drum belonging to the imprinting-unit group.\nA drawback of the aforedescribed system is the low maintenance of register thereof. For example, in the imprinting-unit group, it is barely possible to provide four-color images printed in the satellite printing-unit group with a frame of a special color adjoining the four-color image seamlessly, on the one hand, without any overlap and, on the other hand, without any free space. Between the adjacent colored objects, in the example the four-color image and the frame, small register differences occur, which manifest themselves to the viewer for the most part as disruptive so-called “flashes” of the paper background in the printed product. A significant cause of the register differences is the so-called transfer register of the printing machine system. Because of the necessary clearance and the wear of the connecting element that is caused by the use thereof, accurate-reproduction coupling of the groups is not possible, so that the position of the groups relative to one another at the point of separation turns out slightly differently each time they are coupled. In other words, sheet grippers belonging to the transfer drum have a slightly different position in relation to sheet grippers belonging to the acceptance drum, each time they are coupled. Because of the necessity to be able to couple the satellite printing-unit group to the imprinting-unit group and to uncouple them from one another, in the system described, no adequately good transfer register accuracy, of for example ±10 μm (standard deviation), is ensured, as is usual in the case of other modern sheet-fed offset printing machines with printing units which are coupled rigidly and cannot be uncoupled from one another. In the case of the system described in the foregoing German patent document, the pressman has no possibility of compensating for poor transfer register accuracy.\nFurther prior art is described in U.S. Pat. No. 5,660,108, in the published German Patent Documents DE 43 03 797 A1 (which corresponds to U.S. Pat. No. 5,479,856) and DE 195 03 619 A1."}
{"text": "1. Field of the Invention\nThe invention generally relates to apparatus that may be used in performing joint surgery; with a particular embodiment of the invention being described, without limitation, in the context of knee arthroplasty. More particularly, with respect to knee arthroplasty, the invention relates to apparatus used to properly locate and guide instruments that resect bone in order to achieve a proper cut; and facilitate the proper location and installation of artificial femoral and tibial prosthetic components.\n2. Brief Description of the Prior Art\nTotal knee arthroplasty involves the replacement of portions of the patellar, femur and tibia with artificial components. In particular, a proximal portion of the tibia and a distal portion of the femur are cut away (resected) and replaced with artificial components.\nAs used herein, when referring to bones or other body parts, the term xe2x80x9cproximalxe2x80x9d means closest to the heart and the term xe2x80x9cdistalxe2x80x9d means more distant from the heart. When referring to tools and instruments, the term xe2x80x9cproximalxe2x80x9d means closest to the practitioner and the term xe2x80x9cdistalxe2x80x9d means distant from the practitioner.\nThere are several types of knee prostheses known in the art. One type is sometimes referred to as a xe2x80x9cresurfacing typexe2x80x9d. In these prostheses, the articular surface of the distal femur and proximal tibia are xe2x80x9cresurfacedxe2x80x9d with respective metal and plastic condylar-type articular bearing components.\nThe femoral component is typically a metallic alloy construction (cobalt-chrome alloy or 6A14V titanium alloy) and provides medial and lateral condylar bearing surfaces of multi-radius design of similar shape and geometry as the natural distal femur or femoral-side of the knee joint.\nOne important aspect of these procedures is the correct resection of the distal femur and proximal tibia. These resections must provide planes which are correctly angled in order to properly accept the prosthetic components. In particular, the resection planes must be correctly located relative to three parameters: proximal-distal location, varus-valgus angle and flexion-extension angle.\nU.S. Pat. No. 5,916,219 (hereby incorporated by reference) discloses an apparatus and method for tibial alignment which allows the independent establishment of two separate geometric planes to be used as a reference for the cutting of the tibial plateau during total knee arthroplasty.\nTwo separate frame assemblies with telescoping rods are attached to the tibia with a fixed relative angle between them, thereby allowing alignment with the mechanical axis of the bone. A cutting block is mounted on one of the assembly frames and is positioned against the tibia. Stabilizing pins are then placed in the cutting block, allowing the proper tibial plateau resection plane to be created.\nThe apparatus and method taught in the 219 patent, while solving many prior art problems as indicated therein, has a few disadvantages. The alignment apparatus must be removed prior to performing resection. The device ratchets to discrete locations, preventing a smooth (hereinafter defined as an xe2x80x9cinfinitely adjustablexe2x80x9d) alignment. Further, the device can only be used to resect the tibia and cannot be used for femoral resection.\nRecently, various computerized systems have been introduced to aid the practitioner during different surgical procedures. A typical, commercially available system is described in the attached Appendix. Such systems typically include multiple video cameras which are deployed above and around the surgical site; and a plurality of dynamic reference frame (DRF) devices, also known as trackers, which are attached to body parts and surgical instruments.\nThe trackers are generally LED devices which are visible to the cameras. Using software designed for a particular surgical procedure, a computer receiving input from the cameras guides the placement of surgical instruments.\nThe prior art instruments used for determining the correct planes for tibial and femoral resection in total knee arthroplasty are not well suited for use with computerized systems. The known tools utilize either intra-medullary alignment or extra-medullary alignment techniques and movement in three degrees of freedom is difficult or impossible. Moreover, in order to be useful with computer aided navigation systems, trackers must be attached to the tools. Existing tools do not permit or readily facilitate the attachment of trackers.\nAlthough computer aided navigation systems are superior to unaided visual navigation by the practitioner, computers have known faults. As every computer user knows, any computer can crash or fail in such a way that it may take hours to repair. This is unacceptable during a surgical procedure. Therefore, it is necessary to provide a backup system of some kind so that the procedure may be completed without the failed computer.\nIt is therefore an object of the invention to provide methods, systems and tools for performing femoral and tibial resection and indeed, methods, systems and tools similarly useful in performing surgery on other joints (implicit in the objectives and description of the invention set forth herein, although reference to the femur and tibial resection is specifically made for the sake of illustration).\nIt is also an object of the invention to provide methods, systems and tools for femoral and tibial resection which allow location of a cutting guide relative to three parameters.\nIt is another object of the invention to provide methods, systems and tools for femoral and tibial resection which are infinitely adjustable.\nIt is still another object of the invention to provide methods, systems and tools for femoral and tibial resection which are adapted to be used with computer aided navigation systems.\nIt is also an object of the invention to provide tools which can be used for both femoral and tibial resection.\nIt is another object of the invention to provide methods, systems and tools for femoral and tibial resection which can be used without computer aided navigation systems should such a system fail during surgery.\nIn accord with these objects which will be discussed in detail below, the tools of the present invention (again, described in the context of knee surgery for illustrative purposes only), include anchoring devices for attachment to the femur and the tibia, a three-way alignment guide attachable to the anchoring devices and adjustable relative to three parameters, two embodiments of a resection guide attachable to the alignment guide and equipped with couplings for trackers, a plane probe for examining the resected plane, and apparatus for EM alignment of the resection guide in the event of a computer failure.\nAccording to one aspect of the invention, the tibial anchoring device is specially designed to fit the triangular anatomy of the tibia. According to further aspects of the invention, the two resection guides both include rotatable pin guides to allow selection of anchoring pin location; one resection guide may be slotted and according to another embodiment of the invention, one is not slotted.\nFurther, according to one embodiment of the invention, the apparatus for visual EM alignment includes an EM rod and an EM alignment guide (sometimes referred to as an alignment handle). The EM rod is attachable to the EM alignment guide and the EM alignment guide is attachable to the cutting guide. The EM alignment guide contemplated by a further aspect of the invention has two ends, one of which as adapted for femoral alignment and the other of which is adapted for tibial alignment.\nThe methods of the invention include operating the computer aided navigation apparatus in the conventional manner including attaching one or more trackers to the bone to be resected; choosing a location for the anchoring device with or without guidance from the computer and installing the anchoring device; attaching the three-way alignment guide to the anchoring device; attaching a resection guide to the alignment guide; attaching one or two trackers to the resection guide; locating the resection guide with the aid of the alignment guide and the computer; fixing the resection guide to the bone with pins through the rotatable pin guides; and resecting the bone.\nAfter the bone is resected, the resection plane probe may be attached to a tracker and moved about the resected plane to obtain feedback from the computer navigation system.\nIn the event of computer failure, the methods include attaching the EM alignment guide to the resection guide; attaching the EM rod to the EM alignment guide; and locating the resection guide by visual location of the EM rod rather than by feedback from the computer navigation system.\nIn accordance with the present invention, there is disclosed an apparatus for guiding the resection of a patient\"\"s bone during arthroplasty, the apparatus including a resection guide adapted for guiding a cutting device relative to a patient\"\"s bone during arthroplasty, and an alignment guide coupled to the resection guide and adapted for attachment to the patient\"\"s bone, the alignment guide including a first assembly for positioning the resection guide along a translational path, a second assembly for positioning the resection guide along a first rotational path, and a third assembly for positioning the resection guide along a second rotational path.\nIn accordance with the present invention, there is disclosed an apparatus for guiding the resection of a patient\"\"s bone during arthroplasty, the apparatus including a resection guide adapted for guiding a cutting device relative to a patient\"\"s bone during arthroplasty, and an alignment guide coupled to the resection guide and adapted for attachment to the patient\"\"s bone, the alignment guide including a first assembly for positioning the resection guide along a translational path and a second assembly for positioning the resection guide along a first rotational path and a second rotational path.\nIn accordance with the present invention, there is disclosed an apparatus for guiding the resection of a patient\"\"s bone during arthroplasty, the apparatus including a resection guide adapted for guiding a cutting device relative to a patient\"\"s bone during arthroplasty, and an alignment guide coupled to the resection guide, the alignment guide including a first assembly for positioning the resection guide along a translational path and a second assembly for positioning the resection guide along a first rotational path and along a second rotational path.\nIn accordance with the present invention, there is disclosed an apparatus for guiding the resection of a patient\"\"s bone during arthroplasty, the apparatus including a resection guide adapted for guiding a cutting device relative to a patient\"\"s bone during arthroplasty, and an alignment guide adapted for positioning the resection guide along a translational path and along a plurality of rotational paths."}
{"text": "A. Field of the Invention\nThe present invention relates to an apparatus and method for making a multiple density layers or gradients of fluid in a vessel in a highly reproducible manner using a float that floats on the surface of the fluid within the vessel.\nB. Description of the Related Art\nThere are various fields where it is desirable to have density layers or gradients of fluid within a vessel for such purposes as the separation of matter, determining density, etc. Such density layers include, for example, a solution retained in a vessel where the fluid is divided into a plurality of layers, each layer having differing concentrations of a soluble material or solute. For example, a bottom or first layer of fluid may have a concentration of a solute that is X moles per liter; a second layer immediately above the first layer may have a concentration of 0.8X moles per liter; a third layer above the second layer may have a concentration of 0.6X moles per liter; and a fourth layer having a concentration of 0.4X moles per liter.\nLiquids having gradients of temperature, concentration, density and color have been previously prepared. Liquid density gradients have been used for many years, for a wide variety of purposes, in a number of different industries. The inventor has numerous publications and patents regarding certain aspects of gradient formation and use including Anderson, N. G. Mechanical device for producing density gradients in liquids. Rev. Sci. Instr. 26: 891-892, 1955; Anderson, N. G., Bond, H. E., and Canning, R. E. Analytical techniques for cell fractions. I. Simplified gradient elution programming. Anal. Biochem. 3: 472-478, 1962; Anderson, N. G., and Rutenberg, E. Analytical techniques for cell fractions. A simple gradient-forming apparatus. Anal. Biochem. 21: 259-265, 1967; Candler, E. L., Nunley, C. E., and Anderson, N. G. Analytical techniques for cell fractions. VI. Multiple gradient-distributing rotor (B-XXI). Anal. Biochem. 21: 253-258, 1967.\nA variety of other methods for making density gradients have been developed, and Bock, R. M. and Ling, N.-S., Anal. Chem. 26, 1543, 1954, and Morris, C. J. O. R, and Morris, P., Separation Methods in Biochemistry, Pitman Publishing, 2nd ed. (1976) have reviewed many of these. Only one of these methods allowed gradients to be made from multiple solutions, each having a different combination of reagents (Anderson, et al, xe2x80x9cAnalytical Techniques for Cell Fractions. I. Simplified Gradient Elution Programmingxe2x80x9d, Analytical Biochemistry 3: 472-478, 1962) More recent innovations include the use of pumps and pistons, which are differentially controlled by microprocessors, e.g., the Angelique gradient maker (Large Scale Proteomics Corp. Rockville, Md.). Gradients may also be generated during high speed centrifugation by sedimenting a gradient solute such as cesium chloride or an iodinated x-ray contrast medium such as iodixanol. Gradients may be initially prepared as step gradients and linearized by diffusion, by gentle mixing, or by freezing and thawing. A list of references covering existing methods follows.\nDensity gradients are used to make two basic types of separations. The first separates particles on the basis of sedimentation rate (rate-zonal centrifugation), in which case particles are separated on the basis of the size and density (and to a lesser extent shape) and particles will sediment farther if centrifuged for a longer period of time. The second separates particles on the basis of isopycnic banding density, in which case particles reach their equilibrium density level, and do not sediment farther with continued centrifugation.\nFour types of gradients are in general use with either of these basic methods. The first includes step gradients, made by layering a series of solutions of decreasing density (if the solutions are introduced one above the other), and of increasing density (if the solutions are introduced sequentially to the bottom of the tube). The second type comprises linear continuous gradients usually made by a mechanical gradient maker. These are usually introduced slowly through small tubing to the bottom of the centrifuge tube. Linear gradients for either rate zonal or isopycnic zonal centrifugation are useful for resolving very heterogeneous mixtures of particles.\nThe third type of gradient is non-linear, and may be designed to separate particles having a very wide range of sizes or densities. Non-linear gradient may be designed to separate particles on the basis of both sedimentation rate and isopycnic banding density in the same gradient, in which case some particles reach their isopycnic level at some point in the gradient, while others are still sedimenting. Generally such combined separations involve larger and denser particles which band near the bottom of the gradient, while other smaller, and usually lighter particles are still sedimenting in the upper portion of the gradient.\nThe fourth type of gradient is generated in a high centrifugal field by sedimentation of the major gradient solute, and is usually used for isopycnic banding.\nMany reasons exist for desiring to control gradient shape. Gradient capacity (i.e., the mass of particles which can exist in a zone without causing a density inversion) is a function of gradient slope, and a steep gradient can support a greater mass of particles per unit gradient length than can shallow gradients. The greatest particle mass concentration in a gradient separation usually occurs immediately beneath the sample zone shortly after centrifugation is started. As different particles separate in the length of the gradient, the possibility of an overloaded zone diminishes. For this reason it is desirable to have a short steep gradient section immediately under the sample zone, where the highest gradient capacity is required.\nAn additional reason for desiring to control gradient shape is that when a population of particles is present that differ little in sedimentation rate, these can best be separated by sedimentation through a longer shallower section of the gradient. Such shallow sections are usually near the center of a gradient.\nIn the majority of density gradient separations, the gradients and their chemical composition are designed to optimize the separation of one or a few particles types. This accounts for the very large number of different gradient recipes that have been published for subcellular fractionation. Those used for the isolation of mitochondria, for example, are usually quite different from those used to isolate nuclei. For example, traces of divalent cations are required to control nuclear swelling, whereas such ions are generally deleterious to other subcellular particles. Low concentrations of nonionic detergents remove cytoplasmic contamination from nuclei, but are deleterious to the endoplasmic reticulum. Hence there has been no one procedure or gradient that has been optimized for the systematic separation of the majority of all subcellular particles. There is a need for reproducible means for including in gradients zones containing salts, detergents, enzymes and other reactive substances that would increase the number of different subcellular particles separated in one gradient.\nDensity gradient separations are important in proteomics research. High resolution two-dimensional electrophoresis (2DE) is widely used to produce global maps of the proteins in extracts prepared by solubilizing whole cells or tissues. By careful control of the procedures employed, use of staining procedures which are quantitative, and computerized image analysis and data reduction, quantitative differences in the abundance of individual proteins of xc2x115% has been achieved (Anderson, N. Leigh, Nance, Sharron L., Tollaksen, Sandra L., Giere, Frederic A., and Anderson, Norman G., Quantitative reproducibility of measurements from Coomassie Blue-stained two-dimensional gels: Analysis of mouse liver protein patterns and a comparison of BALB/c and C57 strains. Electrophoresis 6: 592-599, 1985; Anderson, N. Leigh, Hofmann, Jean-Paul, Gemmell, Anne, and Taylor, John, Global approaches to quantitative analysis of gene-expression patterns observed by use of two-dimensional gel electrophoresis. Clin. Chem. 30: 2031-2036, 1984). There is a need for precision subcellular fractionation that will allow changes in abundance of minor proteins to be accurately detected and measured in data which sums the abundance of all proteins found in all of the fractions of one sample.\nThis technology allows changes in gene expression, as reflected in protein abundance, to be studied under a wide range of conditions, and has led to the development of databases of protein abundance changes in response to a wide variety of drugs, toxic agents, disease states. In such studies large sets of data must be acquired and intercompared. Hence all stages in one pharmaceutical study, for example, must be standardized for inter comparability.\n2DE maps of whole cells or tissues typically contain a thousand or more protein spots in sufficient abundance to allow each protein to be analyzed by mass spectrometry and identified and characterized. However, it is known that a very much larger number of proteins are actually present in tissue samples analyzed than are actually observed. The number present varies with cell or tissue type, and is believed to be up to ten or twenty times the number detected.\nDifferent subcellular particles and the soluble fraction of the cell (the cytosol) contain many location-specific proteins which constitute only trace fractions of the total cell protein mass. Hence the total number of proteins resolved from one cell type or tissue could be greatly increased if the 2DE analysis were done on cell fractions rather than on whole cell or tissue extracts as has previously been demonstrated (Anderson, N. L., Giere, F. A., et al, Affects of toxic agents at the protein level: Quantitative measurements of 213 mouse liver proteins following xenobiotic treatment. Fundamental and Appl. Tox. 8: 39-50, 1987). If a drug effect study is to be done on cell fractions, however, the fractionation procedures must be quantitative, in the sense that the same organelles, or even mixtures of organelles are used in all analyses to be intercompared. There exists, therefore, an emerging need for high resolution density gradient separations using precision gradients in proteomics research. Making precision gradients reproducibly and in parallel has proven to be difficult, particularly when the gradients are shallow.\nThe protein composition of tissues such as liver varies diurnally, hence all the tissues from one group of animals are prepared at the same time of day, and, to be comparable, must be fractionated in parallel, on the same time schedule, and, if gradients are to be used, in identical gradients. Further, gradient fraction recovery must also be done from all gradients in parallel, under identical conditions. If the initial separations are done partly or entirely on a sedimentation rate basis, and if the recovered fractions are to then each be isopycnically banded, as is done in two-dimensional or s-xcfx81 fractionation, then these subsequent steps must also be carried out in parallel. This, in turn, requires that the gradients be made in parallel.\nPrecision gradients are difficult to make in practice, and it is further difficult to confirm that a set of gradients are all identical without destroying them for analysis. Existing swinging bucket rotors generally allow six gradients to be centrifuged simultaneously. Larger numbers may be centrifuged if the lower resolution of vertical or near vertical tube rotors is accepted. Therefore if existing density gradient formers are to be used, a set of six or more of them operating in parallel will be required.\nWith any gradient maker, small amounts of turbulence or non-laminar flow typically cause solutions of differing concentrations to at least partially mix, thereby reducing the effectiveness and usefulness of the density layers. There is therefore a need for a method for decelerating fluids flowing into a tube, and for moving them slowly into position to form distinct bands.\nOne of many uses of density layers and gradients is in the fields of cell separation, sub-cellular fractionation and analysis, and density gradient methods are used in molecular biology and in polymer chemistry. Little attention has been paid to forming sets of precision-made gradients that are highly reproducible for cell separation. There is therefore a requirement for precision gradients adapted to cell separation.\nOne high resolution system is disclosed in xe2x80x9cDevelopment of Zonal Centrifugesxe2x80x9d, by N. G. Anderson, National Cancer Inst. Monograph 21, 1966) and employs zonal centrifuge rotors. The rotors are of high capacity, and process one sample at a time. However, the rotor volumes are too high for many applications. Angle head or vertical rotor tubes may also be employed (Sheeler, P., Centrifugation in Biology and Medicine, Wiley Interscience, N.Y., 1981, 269pp) using either step or continuous gradients. However these do not provide the resolution obtained with swinging bucket rotors.\nThere has been no reliable method for reproducibly locating and recovering organelle zones purely on the basis of the physical parameters of sedimentation rate and isopycnic banding density. Mathematical analyses, based on analysis not only of the biological particles separated, but of the gradients themselves have been required. These have been tedious and idiosyncratic to the rotors and conditions employed. The basic problem in preparing density gradients in tubes is that the liquid volume elements of either step (layers), or continuous gradients must be introduced into tubes very slowly or mixing will occur. This problem is only partially overcome by introducing the gradient into a set of tubes in an angle-head rotor during rotation.\nMethods for producing one or a few gradients in parallel have been developed, but fraction recovery is generally done one at a time. The gradients are rarely identical, and it is difficult to introduce the sample layer on top of the gradient without mixing. Hence there is no published data on the quantitative high-resolution protein analysis of cell fractions of animals subject to various experimental treatments. If multiple, parallel identical gradients are to be prepared using gradient engines (for instance, see xe2x80x9cMechanical device for producing density gradients in liquidsxe2x80x9d by N. G. Anderson,. Rev. Sci. Instruments 26: 891-892, 1955) one must have one machine for each tube being filled. Centrifugal gradient distributing heads have been built (see xe2x80x9cA Method For Rapid Fractionation of Particulate Systems by Gradient Differential Centrifugationxe2x80x9d by J. F. Albright, and N. G. Anderson, Exptl. Cell Research 15: 271-281, 1958), however the gradients actually produced tend to be uneven, and a refrigerated centrifuge is required. There is, therefore, a continuing need for simple gradient makers that produce identical gradients in parallel in sufficient number to satisfy current requirements. There is a further need for a simple, disposable and easily sterilizable system for making reproducibly sharp step gradients. An additional need exists for a system or device that can produce very narrow-step density gradients in which diffusion can rapidly and reproducibly even out the steps. A further need exists for a system or device which allows individual gradient steps to be rapidly pipetted into centrifuge tubes, either manually or robotically, and in which the introduced fluid does not disturb the underlying gradient. A still further need exists for a gradient making device in which the composition of the successive layers, while forming a stable density series, differ in composition relative to salts, enzymes, detergents or other reactive materials.\nOne object of the present invention is to provide a rapid, simple and reproducible method and apparatus for forming a multiplicity of liquid density gradients in vessels.\nAnother object of the present invention is to provide a rapid, simple and reproducible method and apparatus for forming a multiplicity of liquid density gradients in vessels for rate-zonal separations, for isopycnic banding separations, or a combination of the two.\nYet another object of the present invention is to provide an apparatus and method for reproducibly producing a plurality of liquid density gradients in a plurality of corresponding vessels, each vessel having a specific predetermined liquid density gradient.\nAn additional object of the present invention is to provide means for making liquid density gradients in which aliquots of a liquid density series are rapidly pipetted into the centrifuge tubes without regard to potential stirring or mixing.\nA further object of the invention is to decelerate the aliquots ejected from pipettes or automatic pipetters, and to cause them to flow evenly into position without disturbing the underlying fluids.\nA further object of the present invention is to provide means for making the linear or complex gradients by making them initially as step gradients having very small density differences per step.\nA further object of the present invention is to produce step gradients in which the steps are so small that diffusion rapidly evens out the gradient.\nA still further object of the present invention is to make the gradient making components disposable and easily sterilizable.\nIt is a further object of the present invention to make possible construction of sets of identical gradients in a short period of time.\nIt is an additional object of the present invention to make possible addition of the sample layer on top of the gradient at any time after the gradient is formed.\nIn accordance with one aspect of the present invention, there is a method for producing liquid density gradients in a vessel using a float within the vessel includes the steps of:\ninserting the float in the vessel;\nintroducing a first liquid into the vessel;\nintroducing a second liquid into the vessel such that the second liquid contacts at least one surface of the float upon entering the vessel, contact between surfaces of the float and the second liquid allowing the second liquid to form a layer above the first liquid thereby forming separate layers of liquid; and\nrepeating the second introducing step with successive introducing steps with a third, fourth and so on liquid.\nThe float used in the above method slows the velocity of fluid such that flow of liquid is laminar thereby limiting mixing of the two liquids.\nIn accordance with another aspect of the present invention, an apparatus for producing liquid density gradients includes a vessel and a float positionable in the vessel. The float is formed with at least one surface that is shaped to inhibit acceleration of fluid introduced into the vessel thereby restricting turbulent flow of the fluid.\nAn outer peripheral surface of the float and the inner surface of the vessel are sized such that in response to fluid being introduced into the vessel above the float, the fluid undergoes capillary action moving downward beneath the float in the vessel."}
{"text": "In many geographical areas, young celery plants are started in seed beds, sowed in closely spaced relationship in the interest of conserving space. Many advantages accompany the use of seed beds and these include the fact that the young celery plants growing from such seeds can easily be shielded from too much sun by the use of certain open fabric material being arrayed in the manner of a canopy from the tops of poles that are spaced around the periphery of the seed beds. Most importantly, proper moisture control may be relatively easily achieved in seed beds, and also it is possible to use much less fertilizer in such an environment than would have been involved had the plants been started in an open field, with the young plants already spaced in the relationship in which they are ultimately intended to grow.\nCelery seeds are quite small, and it has been conventional to build up their size to a standard diameter by a procedure in which controlled amounts of nutrient are applied to the exterior of the seeds. Many machines for dropping the seeds of standardized size at prescribed intervals in seed beds have been proposed, but in all known instances, these machines failed to drop seeds at a consistently even spacing, and furthermore tended to drop the tiny seeds in multiples from time to time. As a result, two, three or more plants often sprouted very close together, resulting in plants that were considerably undersize.\nAs a consequence of the resulting seedlings being of widely varying size, a considerable amount of \"stoop labor\" was always necessitated, for someone must work the seed beds very closely, selecting only seedlings of a uniform size for planting, thus to assure that the final celery crop will be sufficiently uniform in size that it can later be harvested by mechanical means.\nFor many years it has been conventional to then set the selected seedlings out in open fields by the use of large, slow moving machines upon which a number of laborers may work. These laborers sit facing rearwardly, adjacent wheel-like devices that are operationally mounted at above-ground locations, with these devices being caused to rotate in the direction of machine travel, at a speed corresponding to the ground speed of the machine. Various \"fingers\" are mounted around the periphery of each of these wheels, which fingers are arranged during wheel rotation to come close to the ground. The arrangement is such that each laborer, by working rapidly, can place a seedling between each pair of these \"fingers\" as the machine travels through the field. The laborers must be dextrous if they are to avoid creating a gap in the particular furrow they are planting, and, as a matter of fact, it is customary to have one or more persons following behind the machine, placing plants in missed spaces.\nIt is quite obvious that the laborers assigned to place plants in the wheel-like devices have no time to perform any sort of selection with respect to the seedlings they are placing on the wheel-like devices, and therefore in the ground. Accordingly, the laborer responsible for pulling the seedlings out of the seed beds is in effect responsible for the final celery crop being of uniform size. Unfortunately, the plants grown in the seed beds that are passed by as a result of being too large or too small are usually left to wither and die, although this represents a sizable amount of the initial investment of seeds, fertilizer and labor.\nTo all acquainted with this art, it was obvious that what was needed was an automatic seeder able to plant seeds at such a consistant spacing that plants of extremely uniform size would result, thus making it possible for a seed bed to be stripped only once, instead of it being necessary to strip several times, each time making decisions as to which plants were of proper size."}
{"text": "Automotive vehicle disc brake rotors are an example of a manufactured article that may comprise a cast iron (or other durable, friction wear material) annular rotor body attached to a lower density metal hub. Of course, many other manufactured articles comprise, or could comprise, a relatively heavy and durable component that is subjected to loading stress and wear, and a lighter component serving another function.\nIn the case of the vehicle disc brake rotor, a hub is used to attach the rotor body to a vehicle wheel. The hub may be a round cylindrical body (sometimes shaped like a hat) attached to the wheel and aligned with the axis of rotation of the wheel. Attached to the bottom of the hat (for example) and extending radially outwardly is the annular rotor body or disc (brake frictional surface). In operation of the vehicle a hydraulically or electrically actuated brake caliper member, positioned around the outer circumferential edge of the rotor, presses friction pads against opposite sides (cheeks) of the annular rotor disc when it is desired to stop rotation of the wheel. The complimentary hub and rotor body portions of such brake rotors may be cast of a single metal alloy and formed as an integral article. Or the portion subjected to wear may be formed separately of a suitable metal composition and later attached to the complementary portion made of a lighter material. Thus, a complete and integral brake rotor may be formed, for example, of cast iron, or the rotor portion may be made of cast iron and the hub be made of an aluminum or magnesium alloy. An advantage of forming such an article of a relatively heavy and durable portion and a lighter material portion is the overall weight of the article may be reduced.\nA brake disc is subject to frictional heating and to induced vibrations when brake pads are pressed against opposite cheeks of the rotating part. Accordingly, the annular rotor body portion of a disc brake may be formed with radial vents so that cooling air may be pumped through the rotating brake body. And sometimes it is desired to incorporate one or more annular inserts in the annular rotor body to produce columbic friction between contacting surfaces of the embedded insert(s) and surrounding rotor body metal to dampen noise otherwise transmitted by the vibrating brake rotor. These additional structural features of the annular rotor body have significantly complicated the manufacture of low mass articles such as vented and sound dampened brake rotors."}
{"text": "The present invention relates to a new and distinct variety of strawberry (Fragariaxc3x97ananassaDuchesne) plant which is named xe2x80x98Earlibritexe2x80x99 and more particularly to a strawberry plant that is distinguished by its compact plant habit and its high early season (December through February) production of large, bright red fruit when grown in a mild (subtropical) winter climate. Asexual propagation was performed at Dover, Fla. where the selection was made and plants were tested. Contrast is made to xe2x80x98Sweet Charliexe2x80x99 (U.S. Plant Pat. No. 8,729) and xe2x80x98Camarosaxe2x80x99 (U.S. Plant Pat. No. 8,708), standard varieties, for reliable description. This new variety is a promising candidate for commercial success in that it has high early season fruit production like xe2x80x98Sweet Charliexe2x80x99, but tends to have larger and firmer fruit than xe2x80x98Sweet Charliexe2x80x99.\nThis strawberry plant (genotype) originated in a strawberry breeding plot at Dover, Fla. The seed parent was xe2x80x98Rosa Lindaxe2x80x99 (U.S. Plant Pat. No. 9,866), a strawberry variety with a desirable fruit shape and high early season yield potential. The pollen parent was FL 90-38, a non-patented University of Florida breeding selection with an ability to produce attractive fruit early in the season. The seeds resulting from the controlled hybridization were germinated in a greenhouse and the resulting seedlings were planted and allowed to produce daughter plants by asexual propagation (i.e. by runners). Two daughter plants from each seedling were transplanted to raised beds, where they fruited. xe2x80x98Earlibritexe2x80x99 strawberry (as represented by two daughter plants from the original seedling) exhibited large attractive fruit, and therefore was selected for further evaluation. xe2x80x98Earlibritexe2x80x99 was selected from among 207 sibling genotypes as the 100th selection of the 1993-94 season, and thus was designated FL 93-100. It has been asexually propagated by runners, annually, and further test plantings have established that the vegetative and fruit characteristics of the propagules are identical to the initial two daughter plants.\nxe2x80x98Earlbritexe2x80x99, when grown in a subtropical fall and winter climate, is set apart from all other strawberry plants by a combination of three characteristics: compact plant habit, large fruit (averaging greater than 20 gram per berry), and high November through February production (approximately 20,000 pounds of marketable fruit per acre in commercial trials)."}
{"text": "1. Field\nThis disclosure relates generally to semiconductors, and more specifically, to the testing of semiconductors that store data.\n2. Related Art\nAs integrated circuits evolve and incorporate significantly more amounts of circuitry, the ability to reliably manufacture and use reliable fuses is problematic. For example supply voltage values have dramatically been reduced to the point where voltages are not sufficient to reliably blow or open-circuit fuses that rely on an electric current to activate the fuse. An alternative to current or voltage activated fuses is the use of lasers to activate a fuse. However, laser-activated fuses are not an option once an integrated circuit has been packaged. In contrast, fuses that are activated by an operating voltage may be activated after an integrated circuit has been packaged and functionality previously verified. Various advantages exist for waiting until later in the manufacturing process to activate fuses. For example, custom configuration of an integrated circuit may be accomplished wherein certain functions are either enabled or not enabled. Various redundant circuits may be activated. Additionally, fuses may be used by a purchaser of an integrated circuit to implement proprietary security codes in the integrated circuit as well as to create special identification codes for the integrated circuit.\nAdditional problems associated with existing fuse circuits that are used with storage devices include an inability to easily and quickly test the functionality of fuses, both before the fuses are activated and afterwards. Significant circuit area and test time is typically required to check the functionality of programmable fuses."}
{"text": "1. Field of the Invention\nThe present invention relates to an organic light emitting device and a driving method thereof, and more particularly, to an organic light emitting device, which can reduce quality deviation between display devices, and a method of configuring a gamma set of the same.\n2. Description of the Related Art\nIn line with the tendency toward lightweightness and slimness of portable display devices such as notebook computers, cellular phones or portable multimedia players (PMPs) as well as displays for home use, such as a TV or a monitor, a variety of flat display devices are widely used. There are various types of flat display devices, including a liquid crystal display device, an organic light emitting device and an electrophorectic display device. Among the flat display devices, the organic light emitting device is increasingly demanded owing to its various advantages, including low power consumption, high brightness, high contrast ratio, and facilitated enabling of a flexible display.\nThe organic light emitting device implements an image using an organic light emitting diode (OLED) as a light-emitting element. The OLED emits light with brightness corresponding to the current flowing therein. The organic light emitting device includes a plurality of OLEDs and may display an image by controlling gray scales of the respective OLEDs by controlling the current flowing in each OLED. The organic light emitting device may include a plurality of thin film transistors to control the current flowing in each OLED.\nEven if organic light emitting devices are manufactured by the same process, they may have different display quality characteristics in view of brightness and color coordinate system due to differences in OLED characteristics of the respective devices. Accordingly, in order to effectively manage quality of the organic light emitting device, it is necessary to reduce a deviation in display quality."}
{"text": "I. Field of the Invention\nThe present invention relates generally to a shoulder guide for an automotive seat belt.\nII. Description of Related Art\nMany automotive vehicles utilize seat belt shoulder guides for routing the seat belt over the shoulder of the vehicle occupant. These previously known shoulder guides typically include a metal panel welded to a bent wire frame. The panel is then secured to the automotive vehicle frame in any conventional fashion, such as by bolts.\nThese previously known seat belt shoulder guides, however, suffer from a number of disadvantages. One such disadvantage of the previously known seat belt shoulder guides is that such shoulder guides are susceptible to both seat belt rotation within the guide as well as jamming of the seat belt within the guide. Such rotation and/or jamming of the seat belt within the guide is inconvenient for the vehicle occupant in utilizing the seat belt.\nFurthermore, such previously known seat belt shoulder guides were relatively heavy thus adding to the overall weight of the vehicle."}
{"text": "There is a need for lighting apparatus which is low-cost and energy efficient. LEDs (light-emitting diodes) provide light sources which are energy efficient, and advances in LED technology are providing even greater such efficiencies over time. In addition, one advantage of using LEDs as light sources is simply that the smaller packages of such sources translate into much lower material usage within fixtures containing such sources, thus yielding the attendant cost-efficiency. As the technology advances and the cost of LED packages is reduced, such material usage cost benefits will become more important.\nTypical applications for lighting systems are roadway and parking lot lighting in which there are performance requirements that light be as uniformly distributed over areas which are to be lighted while the neighboring regions should be free of light spillage. Thus, along roadways and in parking lots, there is a need to be able to direct light in a laterally-biased direction to avoid so-called trespass lighting on neighboring houses, for example, while providing uniform lighting to the roadways or parking lots.\nVarious LED packages have lenses, some including secondary lenses, for various purposes not related to preferential lateral directing of light. In addition, some limited efforts have been made to develop light-directing LED apparatus using small lenses over small LED light emitter packages, and utilizing lenses intended to redirect some amount of emitted light to one side preferentially. However, such lenses, utilizing lens surfaces having one or more projecting elements and extensive regions of discontinuity, fall well short of requirements to avoid trespass lighting. The measure of trespass lighting includes more than just the amount of light energy falling toward the non-preferential side but also includes how far into the area on the non-preferential that the light falls.\nThe requirements for avoiding trespass lighting may even indicate in some applications that it is preferable to absorb light rather than to allow it to be emitted into the area on the non-preferential side."}
{"text": "The invention relates to an ultrasonic tubular transducer to propagate acoustic waves in a fluid medium, and comprising:                an electromechanical conversion device with active elements associated with at least one metal tube transmitting the vibrations to the fluid medium,        a transmission system of the vibrations between the conversion device and the tube,        and a device for pre-stressing the active elements of said conversion device.        \nTubular ultrasonic transducers enable radial vibrations to be created and waves able to cause cavitation to be propagated in media, in particular fluid media, and find their applications in precision cleaning or sonochemistry, such as extraction of substances or of particle deposition. Their purpose is in particular to generate ultrasounds:                within the fluid and in proximity to the surface to be treated to maximize the efficiency of the installation,        over the whole of the surface to be treated, which may have a length of up to several metres,        with an amplitude that is as uniform as possible for the surface treatment to be homogeneous,        in a direction perpendicular to the surface to be treated to limit losses and to maximize the efficiency.        "}
{"text": "Radio frequency identification (RFID) is a technology that incorporates the use of electromagnetic or electrostatic coupling in the radio frequency (RF) portion of the electromagnetic spectrum to uniquely identify an object, animal, or person. RFID is coming into increasing use in industry as an alternative to the bar code. The advantage of RFID is that it does not require direct contact or line-of-sight scanning. RFID is sometimes also called dedicated short range communication (DSRC).\nIn an RFID system, multiple wireless tags are interrogated by sending information from an interrogating transmitter to the tags and having information transmitted by the tag in response. This is commonly accomplished by having the tag listen for an interrogation and for it to respond with a unique serial number and/or other information. However, it is desirable to extend the range of wireless tags so that it is not necessary to bring each tag close to a reader for reading. Two problems are evident when extending the range of the reading system. One of the problems is that there is limited power available for transmission from the wireless tag, and that if the range is significant, it is possible that many tags will be within range of the interrogating system and their replies may corrupt each other.\nThere are a number of methods and apparatuses that address these problems and enable the fabrication of wireless tags with longer ranges and an ability of the system to handle many tags in the sensitive range.\nIn one application, RFID tags can be used to identify items for purchase, inventory, etc. And, as outlined, there are known methods that enable one specific tag in a group of many tags to be interrogated without corruption by other tags of information sent by that one particular tag to the reader and without accidental transmission of data or commands to other tags of information sent to that particular tag.\nThe least expensive tags usually have read only memory. This is adequate for identifying a tag and for executing a purchase. However, at times, the purchaser might require privacy after the purchase so that another party cannot perform a scan and learn the contents of that purchaser's purse, car or home. Further, the tags may need to be disabled to allow the product to be taken through a security device that recognizes intact tags.\nThus, there are times when it is desirable to disable or destroy an RFID tag after purchase. This disablement may be either partial or complete and may be accomplished by data erasure, modification, etc. However, tags with destructible memory may be susceptible to thievery, in that a thief with the proper equipment could potentially disable a tag. Thus, it is desirable to require some type of validation of authority to disable a tag. Furthermore, when a tag is destroyed it is desirable that other tags within range of the disabling device, typically a reader, are not also destroyed."}
{"text": "Counterfeiting of products, particularly replacement parts for machinery, is more rampant than ever. Verifying the authenticity of a replacement part, or other product, throughout the chain of distribution is therefore more important than ever. Simple printed labels can be easily counterfeited. Therefore, there is a need for sophisticated labels which will foil attempts at counterfeiting. Also, security labels should thwart attempts to tamper with the labels, or to switch them from one part to another."}
{"text": "A common device in the household of today is a shelf unit for holding various items for use in the bathroom, especially in a bathtub or shower area. These units have been manufactured to hold shampoos, various cream and lotion products, wash cloths, shaving products, conditioners, etc.\nOf course, the units must be readily accessible when these items are needed during the bathing process. Consequently, many units have been designed to fit over or attach to the extension arm between the wall and a shower head. A unit of this type is disclosed in U.S. Pat. No. 4,387,811 (1983).\nOther designs have been developed for these shelf units. For example, U.S. Pat. No. 5,921,410 (1999) discloses a collapsible shelf unit which can be attached to a conventional towel rack, hand rail or the like.\nHowever, the prior art units have one or more disadvantages, such as inadequate storage for certain items or lack of structural integrity. In addition, some of the “shower head” units require the user to reach through the shower spray to obtain the desired item.\nThus, there is a need in the industry for an effective and reliable accessory shelf unit for holding various items, especially in a bathtub or shower area."}
{"text": "The present invention relates to a new and distinct cultivar of Sutera plant, botanically known as Sutera cordata, hereinafter referred to by the cultivar name ‘DANCOP20’.\nSutera, of the Scrophulariaceae family, consists of annuals, which are grown in most climate zones in the U.S., and only in zones 9 and 10 if it is a perennial plant.\nThe new Sutera cultivar is a product of a planned breeding program conducted by the inventor, Gabriel Danziger, in Moshav Mishmar Hashiva, Israel. The objective of the breeding program is to develop a new Sutera variety with large, blue flowers and a trailing, compact growth habit.\nThe new Sutera cultivar originated from open pollination in a controlled breeding program by the inventor in 2005, in Moshav Mishmar Hashiva, Israel. Open pollination of the seed parent, an unpatented Sutera cultivar designated ‘CV-05-169’, by an unknown Sutera cultivar, took place in a controlled breeding program in Moshav Mishmar Hashiva, Israel. The new Sutera ‘DANCOP20’ was discovered and selected by the inventor, Gabriel Danziger, as a single flowering plant within the progeny of the stated cross in a controlled environment in Moshav Mishmar Hashiva, Israel.\nAsexual reproduction of the new Sutera cultivar by vegetative tip cuttings was first performed in August of 2005 in Moshav Mishmar Hashiva, Israel, and has demonstrated that the combination of characteristics as herein disclosed for the new cultivar are firmly fixed and retained through successive generations of asexual reproduction. The new cultivar reproduces true to type."}
{"text": "1. Field of the Invention\nOne or more embodiments of the present invention relate to an image forming apparatus and method capable of forming an image on print media.\n2. Description of the Related Art\nIn general, electrophotographic image forming apparatuses such as a laser printer, a digital copier, etc. are apparatuses which emit light onto a photosensitive medium charged to a preset potential, form an electrostatic latent image on an surface thereof, supply the image with a developing agent, develop the image to a visible image, and transfer the image onto a paper to print out the image.\nIn color image forming apparatuses, when images which have different colors are overlapped, if the images are not overlapped at a correct position, edges of the images are blurred, therefore the quality of the images is not good. Since this problem occurs due to various variables such as a changing of a developer, an increasing of the number of copies, a combination thereof, etc., a color registration process is needed to register images in which each color is overlapped and so as to transfer color images onto a correct position.\nFor this type of color registration, light can be emitted with respect to a preset overlap color pattern. Since the degree of reflectivity of light is different according to a color, when light reflected through a photo sensor is received, misalignment of overlapping position of a color pattern can be determined."}
{"text": "Computer devices can enable users to communicate with other users operating other computer devices and to communicate with other computer devices all over the world via, for example, the Internet. Computer devices can use any number of communications technologies to enable network connectivity to the Internet or other networks for desired services, data, media, etc. With the advent of mobile devices and the increasing availability of wireless communication services, users can access the Internet and other networked systems from almost any venue at almost any time.\nEntities such as service providers and enterprises often provide wireless (and wired) services for end users. For example, an enterprise may subscribe to a service provider to receive Internet access for the enterprise's clients. The enterprise may offer unlicensed wireless access (e.g., Wi-Fi) to clients present at a particular location. Some clients who are present at the same location may access the Internet or make calls (e.g., voice, video, etc.) using other technologies, such as licensed cellular access (e.g., 3G, 4G, etc.). Network traffic generated by clients via an enterprise's Wi-Fi offering or other technologies can provide enormous amounts of information associated with the clients, their devices, and the networks being utilized. Harnessing that information and ensuring the security of services that are accessible via the enterprise network, however, presents significant challenges to device manufacturers and network administrators alike."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for manufacturing a highly-crystallized double oxide powder comprising two or more metal elements and/or semi-metal elements, and more particularly relates to a method for manufacturing a highly-crystallized double oxide powder having a uniform particle size and a high purity and composed of a single crystal phase, which is useful as industrial materials in a variety of fields, including various functional materials used in electronics, such as phosphor materials, dielectric materials, magnetic materials, conductor materials, semiconductor materials, superconductor materials, piezo-electric materials, magnetic recording materials, secondary cell-use positive electrode materials, electromagnetic wave absorption materials, catalyst materials, etc.\n2. Description of the Prior Art\nMechanical pulverization has been used in the past to manufacture a double oxide powder. This process involves mixing raw material powders, putting this mixture in a crucible or other such firing vessel and heating at a high temperature for an extended time to bring about a solid phase reaction, and then pulverizing this product in a ball mill or the like. The double oxide powder manufactured by this method, however, is an agglomerate of particles with an irregular particle shape and a broad particle size distribution, and a considerable amount of impurities come from the crucible. Furthermore, the treatment is inefficient because it has to be carried out at a high temperature for a long time in order to raise the homogeneity of the composition. In addition, the particles are often modified on their surface by the mechanical impact and chemical reactions to which they are subjected during the pulverization process, which means that many defects develop at the surface and in the interior of the powder, thereby leading to the lowering of crystallinity and the deterioration of the physical properties inherent in the double oxide.\nA sol-gel process, hydrothermal process, co-precipitation, spray pyrolysis, and so forth are known methods for preventing the generation of a surface modification layer and obtaining a double oxide powder with high crystallinity. Still, a sol-gel process entails high raw material cost because it requires a high-purity raw material, while a hydrothermal process and co-precipitation both take a long time and have low yield per unit of time, so the manufacturing cost is high.\nSpray pyrolysis is a known method for manufacturing a fine oxide powder, in which one or more metal compounds or semi-metal compounds are uniformly dissolved or dispersed in water or an organic solvent, this mixed solution is atomized in the form of fine droplets, these droplets are heated at a temperature higher than the decomposition temperature of the metal compounds and under conditions that will cause precipitation of metal oxides, and the metal compounds are pyrolyzed to produce metal oxide powder. This method makes it easy to obtain a fine spherical powder of uniform particle size, with high crystallinity and no agglomeration, by properly selecting the pyrolysis conditions. Also, since a solution of metal or semi-metal compounds is used as a raw material, advantages are that various metal components can be mixed in the desired ratio and uniformly at the ion level, inclusion of impurities is suppressed, and the particle size of the produced powder can be easily controlled by process control of the spray conditions and so forth. This method is therefore considered a good way to manufacture a double oxide powder, and has been used in the manufacture of a fine double oxide powder for oxide phosphors, oxide dielectrics, oxide magnetic bodies, oxides used for secondary cell positive electrodes, electromagnetic wave absorption materials, catalyst materials, etc.\nNevertheless, when a double oxide powder is manufactured by spray pyrolysis, the retention time of the atomized droplets in a reaction vessel is extremely short, generally lasting only 0.1 to 10 seconds, so in the manufacture of a double oxide with a high melting point, a double oxide with a high sintering reaction temperature, or a double oxide composed of metal elements with low sintering reactivity, the produced particles are not a single phase, and multiple crystal phases tend to be produced, or unreacted material tends to remain. Particularly when many metal compounds are used for the raw material, it is believed that the compounding reaction is impeded because the various components have different pyrolysis temperatures or sintering temperatures. This causes the powder that is produced to have a low crystallinity, or the particle surface may be porous or the particle may be cracked, and it is difficult to fully realize the physical properties inherent to the double oxide, among other problems. One of the problems encountered if the heating is carried out at a higher temperature in an effort to promote pyrolysis and compounding is that the desired crystal phase will not be obtained\nMeanwhile, complex polymerization is a known method for manufacturing a double oxide powder composed of a single crystal phase. With this method, a raw material compound including two or more metal elements that constitute the double oxide, a polyamino chelating agent or a hydroxycarboxylic acid such as citric acid, and a polyol such as ethylene glycol are mixed to produce an aqueous solution, which is heated at a specific temperature of at least 110° C. to bring about an esterification reaction and form a composite polymer of a metal complex in the solution, thereby forming a homogeneous precursor having a uniform element arrangement on the metal ion level. After this, as discussed in Japanese Patent Publication Nos. 6-115934A, 10-330118A, and 10-99694A, for example, powder of the composite polymer of the metal complex is separated, calcined and then fired at a high temperature to obtain a double oxide powder, or as discussed in Japanese Patent Publication No. 10-182161A, this composite polymer is dissolved in water and atomized into a high-temperature atmosphere and pyrolyzed to obtain a double oxide powder. This method, however, is complicated in that it requires a composite polymer of a metal complex to be synthesized prior to the pyrolysis step, and for this polymer to be separated or redissolved. It also requires a large amount of ethylene glycol or other co-precipitant or complexing agent, and the polycondensation reaction takes an extremely long time with a low yield, thereby resulting in a higher manufacturing cost. Furthermore, the homogeneity of the complex ions of the solution may not always be maintained in the removal of the solvent from the produced complex. Particularly when the polycondensation reaction is conducted by gradually removing the solvent over an extended period of time, any gel that is produced may redissolve and form salts with anions present in the solution, and these salts may precipitate and affect the composition, which means that homogeneity is lost.\nIn yet another known method, a double alkoxide, double salt, or the like of the metals that constitute the double oxide is synthesized ahead of time, and this is pyrolyzed, but this method also entails a complex procedure, the metals that can be used must be selected from a narrow range, and it is difficult to achieve homogeneity because of differences in the hydrolysis rates of the respective metal compounds, for example, and therefore a single phase of sufficient crystallinity is not obtained."}
{"text": "1. Technical Field of the Invention\nThe present invention relates generally to sensing motor back EMF signal zero crossing events.\n2. Description of Related Art\nIn the course of operating and controlling brushless DC (BLDC) type motors it is often necessary to monitor the voltage at the motor terminals in order to determine the instant in time that this voltage changes polarity with respect to the circuit common. This event is called the zero crossing. When employing modern control circuitry, this task is often accomplished by connecting these signals to the Analog to Digital Converter (ADC) input pins of a micro processor. The problem which is faced in connecting these signals to the micro processor is that the voltage level of these signals is often quite high, commonly near 350 Volts, while the highest voltage level that can be safely connected to the micro processor is normally in the range of 3 to 5 Volts. It is necessary to protect the control circuitry from damage due to the high voltage while preserving the integrity of the signal so that it can be used for control sensing purposes.\nA method which has been used previously to address this problem is to couple the motor terminal signal to the sensing circuit through a resistive attenuator circuit. An example of such an attenuator circuit is shown in FIG. 1. The resistive attenuator circuit is, in this example, a common resistive divider. One resistive divider is provided for each sensed phase (PHA, PHB and PHC) of the three-phase motor. In a typical case, if the terminal voltage may be as high as 350 Volts and the maximum allowable voltage at the micro processor is only 5 Volts, then the attenuation factor of the circuit must be 70:1 or greater.\nThe foregoing attenuation factor poses a significant limitation because much of the signal strength is lost in the attenuator circuit. At high motor speeds, the BEMF signal present at the terminals is proportionately high so that, even with high attenuation, sufficient signal strength is presented to the control circuitry so that a valid zero crossing detection can be obtained. However, at relatively low motor speeds, the signal which is left after attenuation will not have sufficient strength to obtain a reliable zero crossing detection.\nIf the motor phase signal can be safely coupled to the control circuitry without attenuation then sensing of the BEMF signal can be obtained at considerably lower speeds. This increases the operational speed range of a sensorless BLDC motor drive system and is very desirable.\nAnother technique which has been previously used to couple the terminal voltage to the sensing circuitry in a safe manner is to employ a diode clamping circuit. An example of such a circuit is shown in FIG. 2. The circuit includes a current limiting resistor in series with a Zener diode. One such circuit is provided for each sensed phase (PHA, PHB and PHC) of the three-phase motor. In an alternative implementation, two diodes could be used for connecting to safe voltage reference levels. This circuit does not reduce the base signal level as does the resistive attenuator so operation at relatively low speeds is preserved. Since the sensing circuitry (in this case the micro processor) is only looking for the time of the zero crossing, it is not a concern that the signal is “clamped” to levels no higher than the zener breakdown voltage (normally about 5 Volts) and no more negative than the forward voltage of the diode (about 0.7 Volt).\nThe major problem encountered with the circuit of FIG. 2 is the very high power dissipation encountered in the current limiting resistor. Many micro processors require that the equivalent source resistance of circuits which connect to their ADC inputs should be no higher than 10K ohms in order to get good results. If a value of 10K ohms is used, and the terminal voltage goes to 350 Volts, then the instantaneous power dissipation of the resistor is 11.9 Watts. This power dissipation level is very high and it is not practical to design a circuit in this manner. For motors operating at significantly lower voltages, however, this circuit can be economical and efficient.\nA need exists in the art which addresses the deficiencies of the prior art circuits of FIGS. 1 and 2."}
{"text": "Typically, each battery powered device such as a wireless electronic device requires its own charger and power source, which is usually an alternating current (AC) power outlet. Such a wired configuration becomes unwieldy when many devices need charging.\nApproaches are being developed that use over-the-air or wireless power transmission between a transmitter and a receiver coupled to the electronic device to be charged. Such approaches generally fall into two categories. One is based on the coupling of plane wave radiation (also called far-field radiation) between a transmit antenna and a receive antenna on the device to be charged. The receive antenna collects the radiated power and rectifies it for charging the battery. Antennas are generally of resonant length in order to improve the coupling efficiency. This approach suffers from the fact that the power coupling falls off quickly with distance between the antennas, so charging over reasonable distances (e.g., less than 1 to 2 meters) becomes difficult. Additionally, since the transmitting system radiates plane waves, unintentional radiation can interfere with other systems if not properly controlled through filtering.\nOther approaches to wireless energy transmission techniques are based on inductive coupling between a transmit antenna embedded, for example, in a “charging” mat or surface and a receive antenna (plus a rectifying circuit) embedded in the host electronic device to be charged. This approach has the disadvantage that the spacing between transmit and receive antennas must be very close (e.g., within thousandths of meters). Though this approach does have the capability to simultaneously charge multiple devices in the same area, this area is typically very small and requires the user to accurately locate the devices to a specific area. Therefore, there is a need to provide a wireless charging arrangement that accommodates flexible placement and orientation of transmit and receive antennas.\nWith wireless power transmission there is a need for systems and methods for disposing the transmit antennas in vehicles for convenient and unobtrusive wireless power transmission. There is also a need for adjusting the operating characteristics of the antennas to adapt to different circumstances and optimize power transfer characteristics."}
{"text": "Hepatitis C virus (HCV) infects 1-2% of the population world-wide, forming a chronic infection of the liver (Gower E, et al., 2014, J Hepatol. pii: S0168-8278(14)00526-1). Though largely asymptomatic for much of the infection, HCV-infected persons can develop cirrhosis and liver cancer during the later stages of the infection. The lack of symptoms associated with HCV infections contributes to the fact that the majority of infected persons go undiagnosed. A recent estimate in Canada has predicted that while infection rates will slowly drop over the next 20 years due to HCV screening and the use of anti-viral medications, the number of persons affected by cirrhosis and liver cancer will greatly increase, due largely to the aging of the currently infected population (Myers R P, et al., 2014, Can J Gastroenterol Hepatol. 28(5):243-250). New direct-acting anti-virals (DAAs) will aid in the treatment of HCV, though treatment rates will likely remain low (currently approximately 1.4% of infected persons in Canada, and approximately 5% in Britain and France), treatments are not entirely effective, and DAAs are prone to the development of resistance. The development of a therapeutic vaccine which could be used to treat large at-risk populations would help immensely to stem the oncoming tide of HCV-related disease.\nStudies of HCV-infected chimpanzees and intravenous drug users have indicated that the production of a broad anti-HCV T cell response, potentially in combination with HCV neutralizing antibodies, will eliminate an HCV infection (Ghany M G, & Liang T J, 2013, N Engl J Med. 369(7):679-80). An estimated 20-25% of persons infected with HCV will spontaneously clear the virus, dependent on the production of a suitable immune response. Therefore, there is good reason to believe that a successful therapeutic HCV vaccine can be developed.\nThe largest current hurdle in the development of a successful therapeutic vaccine for HCV is the lack of animal models suitable for testing vaccine efficacy. This stems from the host-restricted nature of HCV infections, being largely restricted to infection in humans and, more specifically, human hepatocytes. Chimpanzees, though infectable, are unlikely to be used to any great extent in coming years, with the United States government no longer supporting chimpanzee use in research, and biopharmaceutical companies discontinuing their use. Vaccine efficacy testing in small animal models, generally mice, has involved the use of HCV recombinant viruses (such as vaccinia virus) (Singh S, et al., 2014, Vaccine. 32(23):2712-21), HCV recombinant lymphoma tumours (Ip P P, et al., 2014, Mol Ther. 22(4):881-90), recombinant expression of HCV antigens in the liver through transgenic modification (Satoi J, et al., 2001, J Virol. 75(24):12121-7), or transient expression construct delivery (Yu W, et al., 2014, Vaccine. 32(27):3409-16). Each of these methodologies has unique advantages and disadvantages; none mimic a chronic HCV infection very well. Recent efforts to generate a transgenic mouse that can be infected with HCV have demonstrated that HCV will only infect and persist in the livers of innate immune-incompetent mice; and even then only at very low levels (Domer M, et al., 2013, Nature. 501(7466):237-41). The use of immune-incompetent animals greatly limits the use of these models in vaccine efficacy testing.\nIt has recently been reported that mouse cells cultured in vitro can be made to be more permissive of HCV replication (Frentzen A, et al., 2014, Hepatology. 59(1):78-88). This was accomplished by over-expressing the liver microRNA, miR-122, in in vitro cultured hepatoma cells that were also deficient for the anti-viral protein, MAVS. The authors suggest that such cells could potentially be used to generate clones that stably replicate HCV, and that ultimately this approach may help to develop an immune-competent small animal model for HCV.\nU.S. Pat. No. 7,416,840 describes cells and cell lines which replicate HCV of non-hepatic human and non-human hepatic origin, and the use of these cells and cell lines to identify anti-HCV agents. The ability of these cells to be used to screen for anti-HCV agents is dependent on the ability of the transfected HCV sequences to self-replicate, as such, the HCV RNA used for transfection of the non-human hepatic cells had already been passaged through human non-hepatic cells and had accumulated permissive mutations allowing it to replicate in these cells.\nU.S. Pat. No. 6,127,116 describes a genetically engineered HCV nucleic acid clone that is capable of replication, expression of functional HCV proteins and infection in vivo and in vitro. The nucleic acid clone includes specifically defined 3′ and 5′ non-translated regions (NTRs) to permit replication of the polyprotein encoding sequences within a host cell.\nInternational Patent Application Publication No. WO 2009/005615 describes a use of the NS4B protein nucleotide binding motif (NBM) of HCV for identifying agents that inhibit a neoplastic cellular phenotype. Expression of NS4B NBM polypeptide in mammalian cells were shown to promote a neoplastic cellular phenotype, thus cells expressing NS4B NBM are described as being useful for in vitro and in vivo methods of screening for anti-cancer agents.\nThis background information is provided for the purpose of making known information believed by the applicant to be of possible relevance to the present invention. No admission is necessarily intended, nor should be construed, that any of the preceding information constitutes prior art against the present invention."}
{"text": "The present invention relates to an image sensor for use in facsimile systems and the like.\nAs an image sensor for use in facsimile systems and the like, an IC type image sensor has been known in which a MOS photodiode array, a CCD image sensor or the like is employed.\nIn addition to such IC type image sensor, there has been proposed a so-called contact type image sensor in which a plurality of photoelectric conversion elements are formed on an insulator substrate and a circuit for scanning these elements is provided either on the same substrate or outside of the substrate. In such contact type image sensor, the width of the array of the photoelectric conversion elements is same as that of a document to be read, and therefore the image sensor of this type can read the document with one picture element to one photoelectric conversion element correspondence by the use of an optical system composed of an optical fiber array, lens array or the like.\nFIG. 1(a) is a circuit diagram of the aforesaid contact type image sensor, FIG. 1(b) is a fragmentary plan view showing the contact type image sensor of FIG. 1(a), and FIG. 1(c) is a sectional view taken along line A--A of FIG. 1(b). In the drawings a photoelectric conversion element 1 is composed of a photoelectric film. This element is equivalently represented by a photodiode PD and a condenser CD as shown in FIG. 1(a). A switching circuit 7 consists of MOS transistors 5 and a shift register 6. These MOS transistors are for switching the photoelectric conversion elements 1, and the shift register 6 is for scanning the MOS transistors 5. Furthermore, reference numeral 2 designates electrodes, 8 signal wire, 9 a load resistor, and 10 a power source for biasing respectively.\nReferring to FIGS. 1(b) and (c), the photoelectric conversion element 1 is constructed in such a manner that a photoconductive film 3 is laminated on the separated electrodes 2 provided on a substrate 11, and an electrode 4 is further superposed thereon, in other words, the photoconductive film 3 is sandwiched between the separated electrodes 2 and the electrode 4.\nThe switching circuit 7 is constructed by integrating the MOS transistors 5 and the shift register 6 onto a semiconductor substrate, and the switching circuit 7 thus constructed is placed on the substrate 11 and connected to the photoelectric conversion elements 1 by lead wire in wire bonding or the like manner.\nIn such construction as mentioned above, when image on a document to be read is focused on the photoelectric conversion elements 1 by means of an optical system (not shown), photoelectric current corresponding to the light intensity is generated in the photodiodes PD, so that electric potentials of the electrodes 2 vary in accordance with the photoelectric current. These electric potentials are transmitted to the signal wire 8 through the MOS transistors which is in ON state, and as the result, they are taken out by the aid of the resistor 9 as picture signals. However, in such contact type image sensor the switching circuit 7 is formed with a substantially same width with that of the array of the photoelectric conversion elements 1. As a consequence, length of the signal wire 8 unfavorably becomes longer so that various noises are induced on the signal wire 8. Such noises principally include noise derived from a clock pulse for driving the shift register 6, spike noise generated at the time when a MOS transistor 5 is turned ON and OFF. These noises are a cause for deteriorating signal to noise ratio of the contact type image sensor and therefore serve as an obstacle to improvement of performance such as speeding up in reading speed of the image sensor."}
{"text": "Unless otherwise indicated herein, the materials described in this section are not prior art to the claims in this application and are not admitted to be prior art by inclusion in this section.\nEfficient management of interacting threads and workers needs fine granularity and awareness of process state and newer processors offer hardware acceleration for thread management. When interacting threads share hardware, various optimizations may be applied to hand off demand between the two. These improvements are typically available if smooth and fast worker handoff can be established between the interacting threads but they allow improvements like pipeline interleaving and fast task switching on successful speculative execution.\nMeanwhile, datacenter multi-worker architectures are currently highly network-centric, designed around multiple minimum size worker/service instances that intercommunicate via messaging or queues so that multiple source workers can assign tasks onto a queue while multiple workers pull tasks off the queue to turn them into the next stage product output. Thus, existing multicore software can use multiple cores well. Existing web services, however, are designed to communicate via messaging capacities like message queues so even if two web services are on neighboring cores they send data between each other via the network—via a trip through at least a virtual router at the Virtual Machine Manager (VMM) level. This approach may be about 4-6 orders of magnitude slower than using intercore hardware."}
{"text": "1. Field of the Invention\nThe invention relates to an apparatus for recording an information signal in a plurality of tracks running in the longitudinal direction of a longitudinal record carrier, such as a magnetic record carrier, to a record carrier and to a method of recording. An apparatus as defined in the foregoing is known from U.S. Pat. No. 4,318,141, document D1 in the list of related documents that can be found at the end of this description.\n2. Description of the Related Art\nConventional organisation of recording of information signals on a longitudinal record carrier is sequential, with information signals being recorded serially in one or more tracks along the length of the record carrier. Information signals are typically recorded in one or more tracks from the beginning of the record carrier to the end of the record carrier. If the information signal to be recorded is longer than the length of the record carrier, the transport mechanism steps the one or more heads to one or more different tracks and the recording process continues back from the end toward the beginning of the record carrier. Consequently, the multiple tracks are recorded in a serpentine like manner on the record carrier. The known apparatus has the disadvantage of a long access time for accessing a particular portions of the information recorded on the record carrier."}
{"text": "This invention relates generally to flip-down type sunglasses which are generally worn in athletic competition. The sunglasses may be moved to a horizontal positon so that they do not obscure the wearer's vision. The glasses may be quickly flipped down into a vertical position to shield the eyes from the sun. While a variety of flip-down constructions have been known in the art, such constructions have been less than completely satisfactory in that the structures required for the mounting and positioning the moveable visor have been difficult to inexpensively mass produce.\nGenerally speaking, in accordance with the invention, an improved construction for flip-down sunglasses is provided. The sunglasses have a fixed part which extends across the wearer's forehead and is attached by the use of earpieces and an elastic strap. Pivotedly mounted to the fixed part is a moveable part carrying a sun visor which may be flipped down by the wearer as desired. The moveable part has a pair of ribs which engage a groove carried on a flexible arm integral with the fixed part. The grooves engage the ribs to position the glasses either horizontally or vertically with respect to the wearer's face.\nAccordingly, it is an object of this invention to provide an improved flip-down sunglass construction that may be inexpensively manufactured.\nIt is another object of this invention to provide an improved flip-down sunglass construction that may be manufactured using molding processes.\nIt is a further object of this invention to provide an improved flip-down sunglass construction that is relatively lightweight.\nStill other objects of this invention will become apparent upon the reading of the detailed specification to follow."}
{"text": "1. Field of the Invention\nThe present invention relates to gambling or wagering games, and more specifically to a card game adapted particularly for casino play or the like. The present game uses two decks of standard playing cards, with players wagering upon the likelihood of a single card matching or being a specific mismatch to any of the cards of two series of cards dealt by the dealer.\n2. Description of the Related Art\nInnumerable games and pastimes have been developed in the past to fill the need for entertainment and relaxation for people. A large percentage of these various games and activities have been developed to provide for wagers to be placed upon the outcome, particularly in the case of various card games which have been developed over the years.\nMost such games are relatively simple, with relatively few rules. This is generally desirable, as it requires less mental effort on the part of the players and also permits the game to proceed at a faster pace, which results in more wagers for the casino or facility operating the game. Examples of such games are blackjack (twenty one), baccarat, and various forms of poker.\nMost such games are comparative, i.e., the suits or values of various hands (dealer, player or players, etc.) are considered against one another, or against some predetermined maximum or minimum value. Poker adds to this by providing a series of different combinations of cards having different ranks relative to one another. As such games have become reasonably well known among those who enjoy such activities, there will occasionally be an attempt to add some further rule or rules to such games, in an effort to increase interest in a well-known game and/or to avoid boredom among players. However, such efforts almost always end up complicating the game to the point that novice players shy away from the game, and experienced players may not wish to play due to the changes in the previously familiar rules and format.\nAnother problem with many wagering card games is that players have only one chance to win in any given hand or round of play. While some players may believe that the payout odds are sufficient to make it worthwhile to place a wager with only one chance of winning per hand, other players may tend to avoid such games in the belief that they have a lesser chance of recovering their wagers. As a result, a number of classic wagering card games, e.g. blackjack and poker, have been modified with the addition of further provisions allowing for the possibility of multiple wagers, and corresponding wins or losses, with each hand or round of play. Of course, this tends to complicate the game, leading to the problems of avoidance by newer players and lack of familiarity by more experienced players, as noted above.\nAccordingly, a need will be seen for a wagering card game which provides a relatively simple and straightforward set of rules, yet which allows players to place multiple wagers upon the outcome of a single hand or round of play, if so desired. The present game responds to this need by providing a wagering game in which a dealer deals two sets or groups of cards from a single deck, and then deals a single reference or “pull” card from a second deck. Players may place wagers upon the likelihood of the pull card matching one of the cards in the first set, or optionally in the second set. Further wagering options allow players to wager upon the likelihood of a specific mismatch to the single card appearing in either the first or second set, i.e., a card having the same value or rank and suit color, but of a different suit, e.g. spades and clubs or hearts and diamonds. Different payouts are provided for different wagers, depending upon the progression of the hand during play.\nA discussion of the related art of which the present inventor is aware, and its differences and distinctions from the present card game invention, is provided below.\nU.S. Pat. No. 4,795,161 issued on Jan. 3, 1989 to Dwight W. Chao, titled “Card Game,” describes a series of numbered cards which may be used as chance means for selecting numbers for play in a lottery game or the like. Some of the cards allow the user to enter one or more preferred numbers thereon. As such, the Chao cards are not strictly a game, but rather allow a person to use the cards for selecting numbers for play in another game.\nU.S. Pat. No. 5,072,946 issued on Dec. 17, 1991 to Mark M. Miller, titled “Method Of Playing A Wagering Casino-Type Card Game,” describes a matching type game in which only two players may play at a time. Players each play a single card, and compare the two. A second card is dealt to the player having the lower valued hand, with the hands being compared again. Play continues until the two hands tie or until one of the hands reaches a predetermined count, in which case the lower valued hand wins. The Miller game teaches away from the present game, in that (1) no wins are provided for a match of any cards, and (2) the game cannot be played by more than two players, making it unsuitable for casino play.\nU.S. Pat. No. 5,566,946 issued on Oct. 22, 1996 to Alan G. Parker, titled “Method Of Playing A Card Game,” describes a modification of a blackjack or twenty one game, in which a player may place two wagers. A conventional ante wager is placed initially, before the cards are dealt. However, players may place additional wagers during play, betting on the likelihood of the count of the players' hands exceeding the count of the dealer's single hand, without exceeding twenty-one. Different payouts are provided, depending upon the difference between the players' and dealer's hands. However, Parker does not provide any payout for any form of match, nor for specific mismatches of a single reference card and one of the cards in the hands, as is the case in the present card game.\nU.S. Pat. No. 5,813,673 issued on Sep. 29, 1998 to Ronald Richardson, titled “Method Of Playing A Card Game,” describes a wagering game in which a single reference card is dealt to the dealer with each player receiving two cards, all from the same deck. The object is for at least one of the player cards to be higher or lower than the single dealer card, as decided upon by the player before the cards are dealt. Richardson does not make any provision for exact matching of rank and suit between the dealer's and players' hands, as only a single card deck is used. Moreover, Richardson does not provide for wagering on the probability of a specific mismatch occurring, i.e., cards of the same rank and suit color but different suits, as is provided optionally in the present game.\nU.S. Pat. No. 5,816,577 issued on Oct. 6, 1998 to Mark Preston et al., titled “Game Of Chance Matching Game And Method Therefor,” describes a game in which each of the players are dealt three cards from one deck, with the dealer dealing a single reference card from a second deck. The object of the game is for players to match the rank or value of the single dealer card. The Preston et al. game differs from the present game in that (1) Preston et al. continue to provide for additional dealer cards, in the event that none of the player cards match during the first round; (2) Preston et al. count any similarly valued cards of different suits as a match; he does not require the suits to match; (3) Preston et al. do not provide for any specific mismatches; and (4) Preston et al. provide an individual hand to each player, rather than dealing collective player hands.\nU.S. Pat. No. 5,839,731 issued on Nov. 24, 1998 to John Feola, titled “Method And Apparatus For Playing A Casino Game,” describes a game in which multiple hands are dealt by the dealer, with the dealer determining specifically which card game (e.g., poker, blackjack, etc.) is to be considered to determine the winning hand of the hands dealt. Players wager upon which of the plurality of hands will win, according to the specific game selected by the dealer. Feola does not provide for any matching of a single dealer reference card and another card from any of the hands dealt, nor does he provide for specific mismatches, as is possible with the present card game.\nU.S. Pat. No. 5,895,047 issued on Apr. 20, 1999 to Paul E. Callahan, titled “Method For A Novel Card Game,” describes a game in which only a single card is dealt to the dealer and to each player, with the high card winning. Ties result in a third card being dealt to the tied player. The Callahan game also provides for simultaneous wins by both the bank (dealer) and player, in the event that a third card is drawn and the third card has a lower rank than the two tied cards of the dealer and player hands. Simultaneous losses are provided if the third card is higher than the tied dealer and player cards. No explanation is provided in the Callahan disclosure, as to how both bank and player can win or lose simultaneously during play of a single card hand. Callahan complicates his game by providing for automatic wins or losses in the respective event that an ace or a two is drawn, unlike the present game.\nU.S. Pat. No. 6,004,205 issued on Dec. 21, 1999 to Antonio Lauretta et al., titled “Match The Dealer,” describes a game in which five cards are dealt to each of the players and to the dealer. The object of the game is for the players to match as many cards as possible between their hands and the dealer's hand. However, the card deck used is non-standard, having a single suit with five sets of ten cards each, from ace through ten. Lauretta et al. modify their wagering provision by allowing payouts for matching cards in a given player's hand, i.e., two, three, four, or five of a kind, unlike the present game. Moreover, Lauretta et al. cannot provide for specific mismatches, i.e., like ranking or value and suit color in different suits, since Lauretta et al. do not provide conventionally suited cards.\nU.S. Pat. No. 6,164,652 issued on Dec. 26, 2000 to Antonio Lauretta et al., titled “Match The Dealer Card Game,” is a divisional patent of the '205 U.S. patent to the same inventors, discussed immediately above. The same points raised in the discussion of the '205 U.S. patent are seen to apply here as well.\nU.S. Pat. No. 6,234,483 issued on May 22, 2001 to Blair Bucan, titled “Method Of Playing A Matching Card Game,” describes a game in which each player controls a single hand of three cards, with the dealer dealing one card at a time from a second deck. The object is for the player(s) to match the cards being turned up by the dealer, with the first player to match the various dealer cards with all three of the cards in his or her hand, winning. The Bucan game differs considerably from the present game, in that (1) each player is responsible for his or her own hand, rather than the dealer dealing a pair of hands for matching cards therein with the single reference card; (2) the Bucan game may prove relatively complex for players, as each player must constantly monitor all three of his or her cards as a series of cards is turned up by the dealer; (3) each hand of the Bucan game would likely extend for some period of time, as a relatively large number of dealer cards would have to be dealt in order for one of the players to achieve a match with all three of his or her cards, particularly if a relatively small number of players are playing; (4) Bucan does not provide any award or payout for a specific mismatch, i.e., cards of like value and suit color but in different suits, as provided by the present game; and (5) Bucan provides for only a single wager before the play of the hand, rather than allowing players to place an additional wager(s) optionally upon other potential outcomes.\nU.S. Pat. No. 6,382,629 issued on May 7, 2002 to Clarence B. Hill, titled “Card Game,” describes a keno-like game in which players do not receive cards, but rather receive scoresheets upon which they mark their anticipated winning numbers. A dealer deals a series of cards face down from a conventional deck (with two jokers) in nine stacks, each containing six cards. The dealer then turns up one card at a time from the stacks, working progressively through each stack. Players win depending upon the anticipated matches they previously marked upon their scoresheets, as in keno.\nFinally, U.S. Pat. No. 6,402,149 issued on Jun. 11, 2002 to Henry H. Chou, titled “Method Of Playing A Magic Seven Card Game,” describes a game using a non-standard deck of cards comprising four suits from ace through six, and two jokers. Each player is dealt two cards, with the object being to attain a hand which is higher, lower, or equal to seven, as predetermined by the players before the hand is dealt. The dealer does not receive any cards in the Chou game, but acts only as a dealer, dealing cards to the players and monitoring the game. Chou does not provide for any matching of cards or specific mismatching, as provided by the present card game.\nNone of the above inventions and patents, taken either singly or in combination, is seen to describe the instant invention as claimed. Thus a card game solving the aforementioned problems is desired."}
{"text": "The present invention relates to a card connector, and more specifically, to a card connector with an ejection mechanism of a push-push type. In the card connector of the invention, it is easy to insert a card.\nPatent Reference has disclosed a conventional card connector with an ejection mechanism of a push-push type. In such a card connector, a tray is disposed to eject outside the card connector for covering a bottom surface and side surfaces of a card, so that the card is inserted into the card connector. When the card is removed from the card connector, the tray is ejected again. Accordingly, a mechanism for ejecting the tray determines how easy a card can be inserted into the card connector.\nPatent Reference: Japanese Patent Publication No. 2000-148927\nIn another type of conventional connector, instead of the tray, an ejector having a plate portion is used. In order to reduce a size of a card connector, the ejector is arranged to slide inside the card connector with a card mounted thereon, thereby not being ejected outside the card connector. Accordingly, in order to easily insert the card into the card connector, it is necessary to provide a large insertion opening for inserting the card.\nIn view of the problems described above, an object of the present invention is to provide a card connector with an ejection mechanism capable of solving the problems in the conventional card connectors.\nFurther objects will be apparent from the following description of the invention."}
{"text": "This invention is relevant to seismic data processing in the field of geophysical exploration for petroleum and minerals. The general seismic prospecting method involves transmission of elastic, or \"seismic,\" waves into the earth and reception of reflected and/or refracted waves at the earth's surface (or, occasionally, in a wellbore) via geophones, hydrophones, or other similar devices (hereinafter referred to collectively as \"geophones\"). The elastic waves may be generated by various types of sources, dynamite and hydraulic vibrators being particularly common. As these waves propagate downward through the earth, portions of their energy are sent back to the earth's surface by the acts of reflection and refraction which occur whenever abrupt changes in impedance are encountered. Since these impedance changes often coincide with sedimentary layer boundaries it is possible to image the layers by appropriate processing of the signals returned to geophone groups.\nMany methods of seismic data processing in use today require calculations using traveltimes to determine information regarding subsurface geology. \"Traveltime\" means generally the amount of time a seismic signal takes to travel from a seismic source to a subsurface reflection point to a seismic receiver. For example, the concept of migration is well known in the art. Simplistically, raw seismic data as recorded are not readily interpretable. While they show existence of formation interfaces, raw data do not accurately inform the interpreter as to the location of these interfaces. Migration, also called imaging, is the repositioning of seismic data so that a more accurate picture of subsurface reflectors is given. In order to perform the migration calculations, the seismic velocities (hereinafter referred to as \"velocities\") of the subsurface at a multitude of points must first be determined or approximated. These velocities are often estimated using traveltime information. Current methods of computing the traveltimes necessary to perform three-dimensional depth migration and associated velocity analyses are inefficient and/or potentially error-prone when applied to the complicated velocity models typically encountered.\nThere are currently two general methods of determining the grid of traveltimes needed to migrate data: the well-known ray tracing methods; and the more recently proposed methods which seek a direct solution of the eikonal equation. As is known to one skilled in the art, the eikonal equation is a form of the wave equation for harmonic waves, valid only where the variation of properties is small within a wavelength, otherwise termed \"the high-frequency condition.\" In rectangular coordinates, the eikonal equation is as follows: EQU (.differential.t/.differential.x).sup.2 +(.differential.t/.differential.y).sup.2 +(.differential.t/.differential.z).sup.2 =S.sup.2 (x,y,z) (1)\nwhere the coordinate axes are x, y, and z; t is the traveltime; and S is the slowness, which is the inverse of velocity. This equation relates the gradient of traveltime to the velocity structure.\nSeismic ray tracing methods are applied to determine traveltimes in most applications used commercially today. Ray tracing equations are linear, ordinary differential equations derived by applying the method of characteristics to the eikonal equation, a technique known to those skilled in the art. Ray tracing allows the determination of arrival times throughout the subsurface, by following raypaths from a source location, which raypaths obey Snell's law throughout the subsurface volume for which the velocity distribution is known. Traveltimes along the rays are then interpolated onto a three dimensional grid of the subsurface.\nThe ray equations may be solved with shooting methods (\"shooting\") or with bending methods (\"bending\"), as are well known to those skilled in the art. Shooting formulates the ray tracing equations into an initial-value problem, where all ray direction and position components are defined at the source location at time zero. Then the equations are recursively solved to trace the rays throughout the medium. Bending is based on Fermat's principle, which states that the seismic raypath between two points is that for which the first-order variation of traveltime with respect to all neighboring paths is zero, and attempts to locate a raypath between two points by determining a stationery traveltime path between them. It formulates the ray tracing equations into a two-point boundary value problem. Shooting is generally more efficient computationally than bending; however, both approaches present difficulties and potential inaccuracies when used to compute the gridded traveltime fields required by three-dimensional depth migration. Three-dimensional depth migration typically requires robust grids of traveltime for high quality images. As used herein, the term \"robust\" means a process which reliably generates accurate grids of traveltimes regardless of velocity model complexity.\nComplications arise in shooting and bending calculations because each ray is computed independently of all others, and because small changes in the angle of incidence may lead to large changes in ray direction. This complication is a manifestation of the nonlinearity of the problem of ray tracing. Thus, with shooting methods, it is difficult to obtain the uniform ray coverage throughout the model that will admit traveltime interpolation onto a grid. Also, rays may cross, indicating the natural multivaluedness of wave propagation. However, when interpolating traveltimes onto a grid, one may not interpolate between different branches of a multivalued wavefront, hence a \"smart\" interpolation, that is, an interpolation geared to adjust its parameters based upon conditions at various locations within the grid, is necessary for a multivalued wavefront. Unfortunately, a \"smart\" interpolation is a difficult task in three dimensions.\nRay shooting itself is quite efficient, but a robust application using it to compute grids of traveltimes may not be efficient. While the bending method will determine a raypath to any point in the grid, it is subject to local, instead of global, minimization, resulting in traveltime errors. Bending is also inefficient for determining three dimensional grids of traveltimes.\nAn alternative to ray tracing is to solve the three dimensional eikonal equation directly, using numerical techniques. This method was proposed in Vidale, J. E., \"Finite-difference Calculation of Traveltimes in Three Dimensions,\" Geophysics, vol. 55, pp. 521-526 (1990), Vidale succeeded in calculating a grid of traveltimes with a finite-differencing scheme based on the rectangular version of the eikonal equation as set forth above in Equation (1), which scheme makes a plane-wave approximation, a simplifying assumption which introduces error into the calculation. An alternative to Vidale's method was proposed in Podvin, P. and LeComte, I., \"Finite-difference Computation of Traveltimes in Very Contrasted Velocity Models: a Massively Parallel Approach and its Associated Tools,\" Geophys. J. Int., vol. 105, pp. 271-284 (1991), which also makes a plane-wave approximation, which approximation is not used in the method of the present invention. This alternative constitutes a finite-difference approximation of Huygens principle. Huygens principle holds that every point on the boundary of a \"cell,\" wherein a plurality of these cells define a subsurface grid, acts as a secondary source emitting an impulse at the moment it is reached by the first wave arrival.\nThese algorithms do not vectorize or parallelize at the loop levels. This means that these algorithms are not set up to contain \"independent\" nested iterations of calculations or \"loops\". A vectorizable algorithm contains single unnested do loops whose iterations may be computed in any order (thus they are independent). A vectorizable and loop parallelizable algorithm contains doubly nested do loops, whose iterations are independent of one another. The inner loop vectorizes and the outer loop parallelizes. Vector loops may be executed very efficiently on any one processor of a vector computer, while parallel loops may be executed efficiently by multiple processors on a parallel computer. A CRAY.TM. Y-MP is an example of such a machine.\nThe two-dimensional graph theory approach of Moser, T. J., \"Shortest Path Calculation of Seismic Rays,\" Geophysics, vol. 56, pp. 59-67 (1991), extends to three dimensions; however, this approach, also a Huygens principle method, is also inefficient.\nVan Trier, J. and Symes, W. W., \"Upwind Finite-Difference Calculation of Traveltimes,\" Geophysics, vol. 55, pp. 521-526 (1991), drew upon fluid dynamics technology and presented a two dimensional upwind (meaning only traveltimes at points in the direction of the negative traveltime gradient, with respect to the current calculation point, are used in the determination of the traveltimes) finite-difference traveltime calculation technique that fully vectorizes, and Popovici , A. M., \"Finite-difference Traveltime Maps,\" Stanford Exploration Project Report vol. 70, pp. 245-256 (1991), presented the extension of the Van Trier and Symes method to three dimensions, while conceding that the algorithm was unstable. The three dimensional algorithm radially extrapolates the three components of the slowness vector (where slowness is the inverse of velocity), i.e., .differential.t/.differential.r, (1/r).differential.t/.differential..theta., and (1/[r sin .theta.]).differential.t/.differential..phi., in a spherical coordinate system, where\nt=traveltime; PA0 r=the radius from the origin point of the spherical coordinate system to the point of interest; PA0 .theta.=the angle that the z-axis makes with the radial line connecting to the point of interest; and PA0 .phi.=the angle from the x axis to a projection of the radial line onto the x-y plane.\nWhile the efficiency of Popovici's algorithm makes it attractive, it is, as Popovici conceded, unstable. It is particularly unstable and inaccurate when applied to the highly variable velocity fields encountered in three-dimensional depth migration applications. Popovici, \"Stability of finite-difference traveltime algorithms,\" Stanford Exploration Project Report , vol. 72, pp. 135-138 (1991), presented a geometrical derivation of the Courant-Friedrichs-Lewy (CFL) stability condition for the two dimensional method. However, no one has been able to improve the Popovici three-dimensional solution so as to remove its inherent instabilities."}
{"text": "Several types of wire based communication networks exist to provide communication among electronic devices. Many of these networks transmit a differential representation of the data over the network. A differential network uses a transmission cable that has a positive and a negative conductor, and positive and inverted representations of the data are sent on the conductors. A differential signal has the advantage of allowing faster data rates because the differential signals traverse lower voltage swings than single ended signals. Also, the data is less susceptible to noise in a differential signal bus because common mode signal noise picked up on the transmission cable is cancelled by sensing only the difference between the positive and negative conductors of the cable.\nOne critical parameter in differential signal wire based networks is the differential cross-over voltage of the signal transmitters. The differential cross-over voltage is the point where the voltage at the output of the positive signal transmitter crosses over with the voltage at the output of the negative signal transmitter. To minimize communication errors from power supply noise, electromagnetic interference (EMI), or signal ringing, the cross-over voltage should be at a point equidistant between the maximum and minimum voltages of the outputs. This point is often referred to as mid-rail.\nIf the network is a wire based serial network, transceivers are used to transmit and receive signals on the same transmission cable. Transmitters of wire based analog transceivers are generally designed with open-loop differential drivers. The drivers are open-loop in that they do not include a feedback mechanism in controlling their output. These transmitters are designed by tuning the cross-over voltage to an optimal mid-rail assuming a nominal process skew and nominal loading on the transmitter outputs. A problem with tuning is that when the transmitter is realized in silicon the cross-over voltage can deviate from the optimal mid-rail value due to undesired process variations or due to asymmetric parasitic off-chip loading. A deviation in the cross-over voltage from the mid-rail voltage value can result in low yield in semiconductor fabrication of the transmitters. A mask iteration may be needed to take into account the non-nominal conditions and to re-tune the cross-over voltage to the mid-rail value.\nWhat is needed is a differential transmitter with a self adjusting cross-over voltage."}
{"text": "This invention relates to airfoil structures for fixed wing aircraft, and more particularly to a leading edge construction for an airfoil.\nPresent day state of the art commercial twin-engined aircraft use leading edge slats for controlling airflow over the fixed wings of the aircraft during takeoff, landing and cruise conditions. However, slats are heavy, expensive, cannot use common parts, and must use the high-speed (i.e., cruise) wingfoil shape in the extended, angled positions at takeoff and landing. Current production Variable Camber Kreuger (VCK) devices are two-position devices (cruise and gapped), which have excessive drag for takeoff. The VCK advantage is a more optimum aerodynamic shape. One specific area of leading edge flap design where even further improvement is desired is in the seal between the panel of the leading edge component and the airfoil, and also between the panel and bullnose of the leading edge component. Often it has been difficult to achieve an excellent seal between the panel/wing and panel/bullnose interface without the use of separate seal elements. An excellent seal formed at these areas, directly between the forward and rearward portions of the panel and the airfoil, would even further enhance the aerodynamic efficiency of the variable camber leading edge when the leading edge flap is in its takeoff and landing positions.\nStill another area where even further improvement would be desirable is in the manner in which the variable camber leading edge flap is deployed from a retracted (i.e., cruise) position into either a takeoff or a landing position. Some designs of variable camber leading edge constructions, when initially moved out of their retracted positions, tend to xe2x80x9cscoopxe2x80x9d air using the undersurface of the panel forming the variable camber leading edge assembly. It would be desirable to provide a leading edge construction which does not operate to scoop air when it is initially deployed from its retracted (i.e., cruise) position into either a takeoff or a landing position.\nThe present invention relates to a Variable Camber Kreuger leading edge assembly for an airfoil of an aircraft. The leading edge assembly includes a variable camber panel which is hingedly secured to a bullnose. The panel/bullnose assembly is suspended from a lower surface of the airfoil by a forward suspension link and an aft suspension link operatively associated with a linkage assembly. The linkage assembly includes connections to the panel at three points along the panel. The linkage assembly is also coupled to the bullnose to cause the bullnose to be folded toward the panel when the leading edge assembly is in a cruise position.\nIn the cruise position the bullnose is folded toward the panel, and the panel forms a generally smooth, continuous lower surface of the airfoil. As the aft suspension link is driven rotationally by a torque shaft of an actuator, the panel rotates outwardly from the lower area of the airfoil such that a leading edge of the panel moves at a greater rate of travel than the trailing edge of the panel. During this portion of movement the bullnose also begins to unfold. The trailing edge of the panel is maintained close to a forwardmost leading edge of the airfoil to thus prevent the leading edge assembly from xe2x80x9cscoopingxe2x80x9d air during this portion of movement.\nThe forward and aft suspension links, together with the linkage assembly, continue to urge the panel/bullnose rotationally, with the bullnose continuing to unfold, until the panel/bullnose is in a landing position. In this position a gap exists between the trailing edge of the panel and the forwardmost leading edge of the airfoil. The bullnose is also rotated during this portion of movement so that it forms a generally smooth, continuous curvature with the panel when the panel reaches the landing position.\nWhen the aft suspension link is rotated further it causes the linkage assembly to urge the panel/bullnose into a xe2x80x9ctake-offxe2x80x9d position. In the take-off position the trailing edge of the panel abuts the forwardmost leading edge of the airfoil to form a smooth seal therebetween. A forwardmost edge of the panel also rests over a trailing edge of the bullnose to also form a smooth seal at this interface. When in the take-off position, the aft suspension link and linkage assembly form an especially strong truss-type load path to provide a highly rigid arrangement that is not susceptible to deflection in response to the high loads experienced during take-off. The sealing action provided by the leading and trailing edges of the panel eliminate the need for separate seals at the panel/bullnose interface and the panel/airfoil interface, thus also serving to reduce the cost of the leading edge assembly."}
{"text": "The present invention relates in the manufacture of electronic vacuum devices to a method of treating fused glass structures, particularly fused glass fiber elements intended for use as a component element of an electronic vacuum device, such as a cathode ray tube or image intensifier tube to provide enhanced physical characteristics. More particularly, my invention relates to the treatment of fiberoptic type faceplates and microchannel plates formed from fused-together fiberoptic glass materials to enhance performance and reliability of the electronic vacuum devices into which they are assembled.\nA fiberoptic faceplate is manufactured form a plurality of individual glass fibers which are stacked together, pressed and heated under pressure to form a uniform structure with a plurality of light transmitting passages so as to be optically transparent between front and back surfaces. Typically, in order to minimize optical \"cross-talk\" between the individual glass fibers in the structure, ensuring that the light entering one end of a fiber exits the opposed end and is not diffused or substantially coupled to other adjacent fibers, a coating of light absorbing material is placed on or adjacent the individual fibers or bundle of fibres prior to fusing them together. This coating usually contains material, such as iron oxide particles and carbon particles, dispersed in an adherent suspension or matrix. the aforedescribed structure is known and is typical of the kinds of faceplates employed as the output display in electronic vacuum tube devices, such as those commonly termed \"night vision devices\" or \"image intensifiers\". As the output display a phosphor coating is applied to one surface of the faceplate, much like the phosphor coated faceplate of a cathode ray tube found in television sets, so as to cause fluorescence under electron bombardment.\nA problem occurs in such a fused fiberoptic structure whenever water vapor or other gases are present in the light absorbing material or are captured in the minute interstices between individual fibers in the plate. Those undesired impurity ingredients create an electrical leakage path between the front and back surfaces of the fiberoptic faceplate, a path through which direct electrical current can flow and a high leakage is undesirable. For example, a low-voltage battery powered power supply is provided in a night vision device to provide the voltages and currents necessary to the operation of the device, the exact circuit details of which are not pertinent to the present invention. Hence, any leakage path between the front and back surfaces of the faceplate is a factor which must be considered by the designer of the power supply for the device since the leakage path poses a power drain on the supply. Not only is battery power in this way wasted, but more importantly it has been found that the leakage resistance may change with time, such as to decrease and drain more current and this detracts from the function of the other elements of the circuit by changing voltages inasmuch as the power supply is intended to supply very minute currents on the order of 200 milliamps at high voltages into a high resistance load. A secondary problem occurs when the extraneous water vapor and active gas molecules migrate from the interstices of the faceplate into the vacuum region of the vacuum device, creating some amount of adverse chemical reaction, termed \"poisoning,\" of electron emitting surfaces, employed as an element therein as is known, thereby reducing device life. From this standpoint the elimination of water vapor and gases from the faceplate prior to assembly into the night vision device is an objective which previously has been sought to be accomplished by the application of high voltage across the plates and under such conditions \"aging\" the plate, the results of which are not to my satisfaction.\nOne object of my invention therefore is to remove water vapor and other gases from the interstices between glass fibers in a fiberoptic plate, as well as to stabilize the leakage resistance of the plate.\nA related component element of a night vision device is the microchannel plate, a known electron \"multiplying\" element which is formed of a fused bundle of hollow glass tubes arranged in a uniform structure defining a surface having on the order of one million minute passages per square inch of surface area. The mircochannel plate is formed originally from a fused bundle of fiberoptic fibers, as was the aforedescribed faceplate. The glass composition of the microchannel plate is such that when the passage walls are chemically treated by prior art processing techniques, including firing in a wet hydrogen atmosphere at high temperatures, the passage walls attain an electron secondary emission ratio greater than unity and a subcutaneous semiconductive layer, which structure is known in the prior art. An unfortunate undesired side-effect of such processing to obtain the desired secondary emission characteristics rendering it useful as an electron multiplying element is that the process creates a diffusion of hydrogen gas molecules into the glass and by-products of water vapor, hydroxyls, carbon monoxide, nitrogen and free hydrogen, all of which are trapped within the bulk of the glass, and are extremely difficult to remove by normal thermal outgassing techniques, such as heating. As a result, the ultimate life of the vacuum device containing the microchannel plate is in part determined by the diffusion or permeation of such trapped gasses to the active photocathode surface causing localized poisoning of the photocathode material.\nAccordingly, another object of my invention is to eliminate extraneous gas and water vapor from a microchannel plate structure."}
{"text": "The present invention relates to a slot machine, and more particularly to a slot machine in which winning is so controlled as to occur with a predetermined probability and in which the number of wins to occur is averaged for any given number of games.\nIn a slot machine, generally a plurality of juxtaposed reels rotating at high speed are sequentially and automatically stopped, or alternatively each of the reels is independently stopped by manipulation of a stop button for each reel. At the time the reels stop, a win decision is made in accordance with the combination of symbols stopped and positioned on a winning line, and coins (including tokens) are paid out in the number that the kind of win indicates.\nIn such a slot machine, the number of all of the possible combinations of symbols is determined by the number and kind of symbols disposed on respective ones of the reels. Among all of the possible combinations, the number of winning symbol combinations determines the probability of getting a win.\nAlthough the overall win probability is so established, it is natural that the probability of getting a win during a given game is governed by temporal chance while repetitively playing game. Thus, despite the overall win probability, the actual win is generated utterly by chance. This means that there is a substantial possibility for the probability of getting a win in any particular game to vary with time. Therefore, if a large number of wins occur during a certain limited number of games, it is known statistically that the probability of getting a win is decreased during the subsequent interval. As a result, a player playing the game during such a subsequent interval may lose interest in the game because of the reduced chance to get a win."}
{"text": "The use of synthetic biomaterials to sustain, augment or completely replace diseased human organs has increased tremendously over the past thirty years. Synthetic implants have cardiovascular applications such as vascular grafts, heart valves, and ventricular assist devices; extracorporeal systems; and a wide range of invasive treatment and diagnostic systems. Unfortunately, existing biomaterials suffer from well-known problems associated with surface-induced thrombosis or clot formation such as thrombotic occlusion and thromboemboli, and infection. Synthetic vascular grafts having a diameter less than 6 mm are currently impracticable, because of potential thrombotic occlusion, and the artificial heart has been plagued with problems of thromboemboli and infection. Advances in the development of artificial organs and artificial vascular grafts have resulted in the need for nonthrombogenic materials.\nThrombosis is initiated by the deposition of a plasma protein layer on the surface of the implanted biomaterial. Thereafter, platelets, fibrin, and possibly leukocytes, adhere to the deposited protein. The interactions between the plasma proteins and the surface of the implant determine the adhesion, the activation and the spreading of platelets, the activation of coagulation, cell attachment and protein deposition. However, at the molecular level, the fundamental forces and interactions of plasma proteins with implants is not well understood.\nThere have been several attempts to create nonthrombogenic surfaces on polymer implants thereby increasing the blood-biocompatibility of implants.\nEarly attempts included precoating the implants with proteins not involved in thrombosis, such as albumin, to mask the thrombogenic surface of the implant. However, such implants loose their nonthrombogenic properties within a short time. Attempts have been made to mask the thrombogenic surface by coating gelatin onto implants such as ventricular assist devices. While the gelatin coating reduced the thrombus formation, it did not adhere to the implant and it did not prevent thromboemboli and infection.\nAttempts have been made to render implants nonthrombogenic by coating the surface of the implant with polyethylene oxide to mask the thrombogenic surface of the implant; it was discovered that such a coating at times also reduced protein adsorption. While this reduced thrombogenesis, the coupling of polyethylene oxide to the surface of the implant involves very complex procedures, and the coated implants do not consistently exhibit protein resistance.\nThere have been many attempts to prepare nonthrombogenic surfaces by attaching heparin to biomaterials, because of heparin's potent anticoagulant properties. However, each method requires that the implant surface be first modified by attachment of a coupling molecule before heparin can be attached. For example, the positively charged coupling agent tridodecylmethylammonium chloride, is coated onto an implant, which provides a positively charged surface and allows heparin which has a high negative charge density, to be attached. However, the heparin slowly dissociates from the surface, to expose the positively charged, TDMAC surface which is particularly thrombogenic. The TDMAC attracts platelets and other cells; cells surfaces have a high negative charge density. Thus the TDMAC heparin coated implant is successful only for short term implants such as catheters.\nImplants coated with heparin coupled to coupling molecules typically have limited anti-thrombogenic effectiveness because commercial heparin preparations contain the protein core and because many heparin molecules which having no anticoagulant activity. As a result, the surfaces soon become covered by adsorbing protein on exposure to blood, thus neutralizing the anticoagulant activity of the active heparin molecule.\nIt is desirable to have implants which resist plasma protein deposition, and to have a simple procedure for modifying the surface of implants. Nonthrombogenic implants would reduce the need for aggressive anticoagulant therapy, improve the performance of implants, particularly cardiovascular prosthetic devices, and encourage the development of devices not currently feasible."}
{"text": "In order to effectively utilize radio wave resources in a mobile communication system, it is required to compress speech signals at a low bit rate. On the other hand, it is expected from the user to improve quality of communication speech and implement communication services with high presence. In order to implement this, it is preferable not only to improve quality of speech signals, but also to be capable of encoding signals other than speech, such as audio signals having a wider band with high quality.\nFurther, in an environment where various types of networks are present, a speech coding scheme is required that can flexibly support communication between different networks, communication between terminals utilizing different services, communication between terminals having different processing performance, and conversational communication at multipoints as well as communication between two parties.\nMoreover, a speech coding scheme is required to be robust against transmission path errors (in particular, packet loss in packet switching networks typified by IP networks).\nOne speech coding scheme satisfying such requirements is the bandwidth scalable speech coding scheme. The bandwidth scalable coding scheme is a coding scheme that encodes speech signals in a layered way, and a coding scheme where coding quality increases in accordance with an increase in the number of coding layers. The bit rate can be set variable by increasing or decreasing the number of coding layers, so that it is possible to effectively use transmission path capacity.\nFurther, with the bandwidth scalable speech coding scheme, it is only necessary to receive at least the data coded by a base layer at a decoder side, and it is possible to allow to some extent information coded by additional layers being lost on the transmission path, and therefore the bandwidth scalable speech coding scheme provides robustness against transmission path errors. Further, the frequency bandwidth of speech signals to be encoded also becomes wider in accordance with an increase in the number of coding layers. For example, for a base layer (i.e. core layer), a coding scheme for telephone band speech of the related art is used. Further, in additional layers (i.e. enhancement layers), layers are configured so that wideband speech which has a bandwidth such as 7 kHz can be encoded.\nIn this way, with the band scalable speech coding scheme, telephone band speech signals are encoded in the core layer, and high-quality wideband signals are encoded in the enhancement layers, so that it is possible to utilize the bandwidth scalable speech coding scheme for both telephone band speech service terminals and high-quality wideband speech service terminals and support multipoint communication including the two kinds of terminals. Further, the coded information is layered, so that it is possible to increase error robustness by devising a transmission method, and readily control the bit rate on the encoding side or on the transmission path. Therefore, the bandwidth scalable speech coding scheme draws attention as a speech coding scheme for future communication.\nThe method disclosed in non-patent document 1 is given as an example of the bandwidth scalable speech coding scheme described above.\nIn the bandwidth scalable speech coding scheme disclosed in non-patent document 1, MDCT coefficients are encoded using a scale factor and fine structure information for each band. The scale factor is Huffman encoded, and the fine structure is subjected to vector quantization. Auditory weighting of each band is calculated using a scale factor decoding result, and the bit allocation to each band is decided. The bandwidth of each band is non-uniform and set in advance so as to become wider for a higher band.\nFurther, transmission information is classified into four groups as described below.\nA: Core codec coding information\nB: High-band scale factor coding information\nC: Low-band scale factor coding information\nD: Spectrum fine structure coding information\nFurther, the following processing is carried out on the decoding side.\n<Case 1> When information for A cannot be received completely, decoded speech is generated by carrying out frame erasure concealment processing.\n<Case 2> When only information for A is received, a decoded signal for the core codec is outputted.\n<Case 3> When information for B is received in addition to the information for A, a high band is generated by mirroring the decoded signal for the core codec and a decoded signal having a wider bandwidth than the decoded signal of the core codec is generated. Decoded information for B is used in generation of high band spectrum shapes. Mirroring is carried out at a voiced frame, and is carried out so that the harmonic structure does not collapse. The high band is generated at an unvoiced frame using random noise.<Case 4> When information for C is received in addition to information for A and B, the same decoding processing as in case 3 is carried out using only information for A and B.<Case 5> When information for D is received in addition to the information for A, B and C, complete decoding processing is carried out at bands where all information for A to D is received, and a fine spectrum is decoded by mirroring a decoded signal spectrum on the low band side at bands where information for D is not received. Even if the information for D is not received, it is possible to receive the information for B and C, and this information for B and C is utilized in decoding of spectrum envelope information. Mirroring is carried out at a voiced frame, and is carried out so that the harmonic structure does not collapse. The high band is generated at an unvoiced frame using random noise.Non-patent document 1: B. Kovesi et al, “A scalable speech and audio coding scheme with continuous bit rate flexibility,” in proc. IEEE ICASSP2004, pp. I-273--I-276."}
{"text": "Automated keyword generation provides valuable advantages when producing searchable content, such as advertising and media content. When a user is browsing on the internet, the usual method by which a user reaches a website is through a search engine, and one factor in how prominently the search engine displays the website is the quality of keywords associated with the website. For search engine marketing (SEM) campaigns, bidding on the right keywords is extremely important. While keyword suggestion tools currently exist, these existing tools often produce overly broad or inapplicable keyword results. Further, in the case that a content provider is generating large amounts of searchable content, any amount of time refining keyword results from existing keyword suggestion tools quickly becomes a major hurdle. Specifically, within the context of generating automated advertising campaigns for quickly changing sources of content, it becomes extremely advantageous for a user to minimize time required to focus and refine keyword generation results to any content produced or content of interest."}
{"text": "1. Field of the Invention\nThis invention relates to an outboard motor control apparatus, particularly to an apparatus for controlling an outboard motor having a torque converter.\n2. Description of the Related Art\nIn recent years, there is proposed an outboard motor having a torque converter interposed between an internal combustion engine and drive shaft to amplify output torque of the engine and then transmit it to the drive shaft for enhancing acceleration performance, etc., as taught, for example, by Japanese Laid-Open Patent Application No. 2007-315498 ('498). In this conventional technique, the torque converter includes a lockup clutch."}
{"text": "Fluorocarbon based fluids, including chlorofluorocarbons (“CFCs”), hydrochlorofluorocarbons (“HCFCs”), and hydrofluorolefins (“HFOs”), have properties that are desirable in industrial refrigerants, blowing agents, heat transfer media, solvents, gaseous dielectrics, and other applications. For these applications, the use of single component fluids or azeotrope-like mixtures, i.e., those which do not substantially fractionate on boiling and evaporation, are particularly desirable. It is also considered important in many applications, including with respect to heat transfer fluids, blowing agents, propellants, solvents and aerosols, that any potential substitute also preferably possess those properties present in many of the most widely used fluids, such as excellent functional properties (for example, heat transfer properties in the case of heat transfer compositions), chemical stability, low- or no-toxicity, low- or no-flammability and/or lubricant compatibility, among others.\nUnfortunately, suspected environmental problems, such as global warming and ozone depletion, and other potential problems such as a flammability level that is higher than desired, have been attributed to the use of some of these fluids, thereby limiting their contemporary use. Hydrofluoroolefins (“HFOs”) have been proposed as possible replacements for such CFCs, HCFCs, and HFCs. However, the identification of new, environmentally-safe, non-fractionating mixtures comprising HFOs are complicated due to the fact that azeotrope formation is not readily predictable. Therefore, industry is continually seeking new HFO-based mixtures that are acceptable and environmentally safer substitutes for CFCs, HCFCs, HFCs and certain HFOs and mixtures of these.\nThis invention satisfies one or more of the above-noted or other needs."}
{"text": "Computer systems generally include one or more processors interfaced to a temporary data-storage device such as a memory device and one or more persistent data-storage devices such as disk drives. Each disk drive generally has an associated disk controller. Data is transferred between the disk drives and the disk controllers. Data is also transferred between the disk controller(s) and the memory device over a communications bus or similar.\nData organization in a computer system such as that above is important in relational database systems that deal with complex queries against large volumes of data. Relational database systems allow data to be stored in tables that are organized as both a set of columns and a set of rows. Standard commands are used to define the columns and rows of tables and data is subsequently entered in accordance with the defined structure.\nThe defined table structure is locally maintained but may not correspond to the physical organization of the data. In a parallel shared nothing relational database data can be stored across multiple data-storage facilities, each data-storage facility in turn including one or more disk drives. Data partitioning can be performed in order to enhance parallel processing across multiple data-storage facilities. The intent behind data partitioning is to evenly distribute the data among all computational elements such that performance scales linearly as more computational elements are added. The database could be hash partitioned, range partitioned, round-robin partitioned or not partitioned at all.\nHash partitioning is a partitioning scheme in which a predefined hash function and map is used to assign rows in a table to respective processing modules and data-storage facilities. The hashing function generates a hash or partition bucket number and the partition numbers are mapped to data-storage facilities. Range partitioning is a partitioning scheme in which each data-storage facility manages the records falling within a range of values. Round Robin partitioned is a partitioning scheme where the data-storage facility is picked in a round robin fashion. No partitioning means that a single data-storage facility manages all of the rows.\nOne drawback of current systems is that the mapping of partitions to the data-storage facilities on which the rows are stored is often required to be physically defined in advance especially when the mapping is not to all the data-storage facilities but to a subset of them. The mapping of partitions to data-storage facilities is therefore static. A static mapping to specific user defined data-storage facilities often leads to uneven distribution of rows over the data-storage facilities. This in turn has the potential to increase the execution time for complex queries.\nUsers of relational database systems require the minimum time possible for execution of complex queries against large amounts of data. In a parallel shared nothing relational database system it is often important to evenly allocate both table rows and free space across multiple data-storage facilities."}
{"text": "In view of the problems associated with prior techniques such as dry or wet pressing and slip casting, thermoplastic injection molding has been increasingly used to form ceramic articles. Greenbodies prepared in this manner show much improved green strength. The process involves forming a ceramic greenbody by injecting into a mold a mixture of ceramic powder dispersed in a thermoplastic polymer that is held at a temperature above the softening point of the polymer. The mold is held at a temperature below the softening point of the polymer to solidify the mixture of powder and polymer in the mold. A typical ceramic powder/thermoplastic polymer mixture has a very high content of ceramic particles, typically from about 50 to about 87 volume %, and a minimum of the binder material.\nThe part is then removed from the mold and placed in a furnace to burn out the binder. The molded article is heated very slowly through the decomposition temperature of the polymer, e.g., at 10.degree. C./hour. The very slow heating is required to prevent deformation or \"slumping\" of the molded article once the furnace temperature exceeds the softening temperature of the polymer. Often, external physical support is needed to prevent slumping during the binder burnout step. Following the binder burnout step, the resulting porous greenbody is sintered, with concurrent shrinkage, to a dense ceramic part with the same shape as the molded object.\nThe use of an organic thermosetting resin as the binder for the ceramic particles has also been suggested. For example, U.S. Pat. No. 4,627,945 discloses injection molding of compositions that are cured in the mold by the addition of a curing agent to a mixture of ceramic powder and thermosetting phenolic resin binder. Thermoset molding has an advantage over conventional thermoplastic molding in that the greenbody is not prone to \"slump\" during the sintering process, since the thermoset binder, once cured, no longer has a softening point. The thermoset organic binder must, however, still be completely burned out of the molded article prior to the sintering step.\nReaction injection molding (RIM) has also been adapted for forming shaped ceramic greenbodies. U.S. Pat. No. 4,906,424 discloses a RIM process for molding a mix of ceramic powder and a polymerizable, low viscosity, multifunctional organic acrylate monomer or mixtures of monomers. The ceramic-monomer mixes are formulated to be highly filled, i.e., greater than 50 vol. %, with ceramic powder, yet have adequate fluidity to be processed at ambient temperature and readily injected into a hot mold. Useful monomers are those that are liquid at room temperature and can be polymerized to irreversibly solidify the fluid composition in the mold. The part is then ejected from the mold and subjected to subsequent post-curing, binder removal, sintering and, if needed, machining to produce a dense ceramic part.\nHowever, organic binders such as polyacrylates must be burned out of the molded part in the process of converting the part to a dense, sintered ceramic article. The carbon-containing char that would otherwise remain in the sintered body would have a deleterious effect on the structural integrity and high temperature performance of the sintered part. Often, the carbon in the binders previously disclosed for RIM processes cannot be completely eliminated in the firing step. In addition, removal of an organic binder can cause structural defects in a sintered part due to voids formed from the rapid generation of volatile materials in the binder burnout step. A further complication arises in fabricating sintered parts of well-defined dimensions. Excessive shrinkage occurs when a high fraction of a ceramic greenbody must be removed in a binder burnout step. When the part finally densifies at high temperatures, dimensional distortion can be extreme, requiring a complex mold design.\nBinder systems that contribute to the ceramic body (\"non-fugitive\" binders) have been used in traditional molding methods, although not in RIM processes. For example, U.S. Pat. Nos. 4,689,252; 4,722,988 and 4,772,494 disclose a crosslinkable silazane polymer that can be cured and subsequently pyrolyzed to convert the polysilazane to a silicon nitride-containing ceramic material. The silazane polymer can be used for coating or impregnating a substrate, making silicon nitride-containing ceramic fibers or as a sinterable binder for ceramic powders.\nThe use of silicon nitride ceramics in a number of high temperature structural applications has been proposed. The advantages of this material in such applications include its higher relative flexural strength and fracture toughness at elevated temperatures. Unfortunately, since silicon nitride is a mostly covalently bonded ceramic, it is difficult to densify fully in its pure state, regardless of whether the unsintered silicon nitride is in the form of a powder compact or the char that is formed by pyrolysis of a silicon nitride ceramic precursor. Additives are necessary to promote a glassy silicate grain boundary phase that aids in densification. It is the presence of the glassy silicate phase that limits the performance of silicon nitride at high temperatures. This glassy phase softens and melts with catastrophic effects on the mechanical properties of the ceramic.\nOne method for eliminating this glassy phase is described in U.S. Pat. No. 4,264,550, where a mixture of silicon nitride powder containing SiO.sub.2 as an oxide surface coating and Y.sub.2 O.sub.3 powder is heated to 1000.degree. to 1400.degree. C. under a pressure of at least 2000 psi to permit a nucleating reaction to take place. The mixture is then heated to a temperature of 1680.degree. to 1750.degree. C. under pressure. The resulting pressed body is claimed to contain fully crystallized grain boundary phases of Si.sub.3 N.sub.4, SiO.sub.2, Y.sub.2 O.sub.3.\nVarious silicon nitride compositions containing metal silicides have been disclosed. For example, U.S. Pat. No. 5 4,407,971 discloses a sintered ceramic body comprising yttrium oxide, aluminum oxide, aluminum nitride and 0.1 to 5% by weight of at least one silicide of Mg, Ca, Ti, V, Cr, Mn, Zr, Nb, Mo, Ta and W. U.S. Pat. No. 4,612,296 discloses a high toughness silicon nitride sintered body containing at least one silicide or carbide of Ti, Zr, Hf, V, Nb, Ta, Cr, Mo and W. U.S. Pat. No. 4,983,554 discloses a sintered silicon nitride ceramic containing aluminum nitride, yttrium oxide, and molybdenum disilicide. U.S. Pat. No. 5,023,214 discloses a sintered Si.sub.3 N.sub.4 product containing a sintering aid and 1 to 80% of a silicide of Fe, Ni or Co. U.S. Ser. No. 07/592,713 filed Oct. 4, 1990, discloses a process for preparing a sintered silicon nitride ceramic containing 2% to 6% of a silicide of Fe, Ni or Co, at least 50% of which is a high metal content silicide.\nThe prior art does not teach a method for (1) rapidly injection molding a high solids, nondilatant dispersion of silicon nitride, a silicate glass-forming sintering aid, and a high metal content transition metal silicide in a curable, liquid silicon nitride ceramic precursor binder at a low temperature, (2) subsequently curing the precursor and (3) sintering the molded article with concomitant conversion of the ceramic precursor binder to a ceramic containing crystallized grain boundary phases."}
{"text": "Such an optical system is described in German Published Patent Application No. 198 08 273. With this known measuring system, having as its essential component a fiber optic short coherent multi-wavelength heterodyne interferometer, a measuring probe directed at the object surface, and a reference probe to be directed at a plane table on which the component having the surface is mounted, are connected to a beam delivery section having a short coherent broadband light source and a modulation interferometer receiving the light therefrom. Position relationships and changes between the object surface and the bearing device are detectable with this probe system of measuring probe and reference probe and may be taken into account during the surface measuring procedures to reduce and/or rule out measuring errors resulting from such deviations in position. However, other relative local changes between the outputs of the probe system points on the surface, e.g., inaccuracies in the probe guidance or relative movements between probe and component in various degrees of freedom during the measuring procedure, may have effects on measurement of the surface geometry of a component or multiple components in a fixed positional relationship to one another, e.g., distance or roughness or roundness. Such influencing parameters on the measuring procedures are difficult to rule out, in particular in measurements of shape and/or distance related to the manufacturing process.\nToday mainly pneumatic measurement methods are used for a robust measurement of shape and distance and/or diameter of guide bores, for example, that will have relevance for the manufacturing process. However, the precision of these methods is technically limited and will probably not be able to comply with increased tolerance demands in the future. However, for high-precision measurements, high-resolution measurement methods in the nanometer range are needed, but such technology may currently be installed only in precision measuring laboratories having vibration isolation.\nMainly punctiform distance sensors are used for the high-resolution determination of shape and diameter on guide bores; the measuring heads of these distance sensors are capable of very accurate determination of the distance from the surface of a component to a reference point. Such punctiform distance sensors must essentially be combined with a shape-measuring machine which guides the sensor mechanically along a defined path over the surface of the component to be scanned to measure the shape of the component. The accuracy of such measurement methods is based essentially on the precision in guidance and the mechanical stability of the shape-measuring machines.\nOptical sensors, in particular those that operate by an interferometric principle, have recently made advances as punctiform distance sensors replacing tactile scanners because of the non-contact measurement principle and the resulting advantages. Of the interferometric sensors, in particular the multi-wavelength heterodyne interferometer having a short coherent light source is a rapid punctiform measurement method that has already been tested industrially and is capable of measuring rough industrial surfaces. However, confocal measurement methods are also known as optical measuring systems."}
{"text": "1. Field of the Invention\nThis invention relates generally to a device and a method for measuring optically the physical properties of solids, liquids and gases, and more particularly to a device and a method for measuring optically the characteristics of a thin film spread on a liquid surface. The present invention is utilized for measurement of the light absorption characteristics which are the basis for various characteristics analysis of thin films, for example, for characteristic analysis of a monomolecular film spread on a liquid surface to be built up during formation of a monomolecular built-up film.\n2. Related Background Art\nAs a device for measuring light absorption characteristics of a material, there is a device for determining light absorption characteristics from transmittance and reflectance. However, when light is irradiated on a material to be measured, light scattering occurs in addition to transmission and reflection of light, and thus direct measurement of the absorbed light component is important in evaluation of light absorption characteristics to achieve a higher precision of measurement.\nAs a device for measuring directly the light absorbing components, there is Photoacousatic Spectroscopy (PAS) or Photothermal Radiometry (PTR), which is a measuring device utilizing intermittent conversion of light energy to thermal energy according to a radiationless relaxation process of the light energy absorbed by the material to be measured when light is intermittently irradiated thereon.\nThere is also a device called Photothermal Deflection Spectroscopy (PDS) for measuring directly the light absorbing components. This PDS device works based on the phenomenon that refractive index changes by temperature distribution caused within the material to be measured or in the vicinity thereof by heat generated by light absorption in the material, whereby the light introduced thereto is deflected. That is, on the measuring site of a material, exciting light which changes the refractive index by creating a temperature distribution by heat generation when the light is absorbed and a probe light for measuring the degree of deflection in terms of the change in a refractive index are irradiated, and the light absorption characteristics of the material to be measured are measured from the wavelength of the exciting light and the deflection of the probe light. The device allows to provide a material to be measured and a detecting system independently of each other and is suitable for measurement at the site where the material is prepared and for remote measurement apart from the site. The basic principle of the present invention is common to this PDS device.\nThe above-mentioned PDS device is inclusive of the two types of the transverse type and the collinear type according to arrangement of the exciting light and the probe light. Both of the types measure the deflection degree of the probe light corresponding to the amount of exciting light absorbed by the material to be measured as mentioned above, and a position sensitive detector (PSD) is frequently used as the detector.\nFIG. 37A shows an example of the collinear type, in which the exciting light 2 emitting from the exciting light source 3 is made intermittent or accentuated with the light intensity modulator 4, condensed at the lens 7b to irradiate the material to be measured 1. The probe light 5 emitting from the probe light source 6 transmits through the measuring site of the material 1 on which the exciting light 2 is irradiated by means of an optical path controller such as the lens 7a or a mirror, etc. to reach the detector 8, where the deflection degree when deflected as shown by the dotted line is measured. FIG. 37B is an example of the transverse type, and the probe light 5 is irradiated in parallel to the surface of the material to be measured 1, as different from the collinear type, otherwise being the same as the collinear type.\nThe PDS device can be theoretically dealt with by solving the thermal conduction equation within the material to be measured. The deflection degree measured in terms of the deflection angle .phi. is proportional to the exciting light intensity, the temperature coefficient of the refractive index (.differential.n/.differential.n), and the temperature gradient (.differential.T/.differential.x) in the region where the probe light passes, etc. The item proportional to the light absorption coefficient of the material to be measured is included in (.differential.T/.differential.x). Also, the (.differential.n/.differential.T) can take either a positive or negative value depending on the material to be measured and this means that the deflection angle can be either positive or negative.\nFIG. 38 is a longitudinal sectional view showing a construction example of one-dimensional PSD. In FIG. 38, the one-dimensional PSD constitutes a uniform resistance layer 74 of P layer on the surface of a flat plate silicon 73, and has electrodes X.sub.1 and X.sub.2 provided on both sides thereof, having also a common electrode 76 at the N layer 75 on the back.\nFIG. 39 is a schematic illustration showing its actuation principle. The light-forming electric charges made correspondent to the incident position of the light Q reaches the above resistance layer 74 as the photocurrent corresponding to its energy to be divided in inverse proportion to the distances from its position Q to the take-out electrodes X.sub.1 and X.sub.2 at the both ends and outputted from the respective electrodes. If the photocurrent is defined as I.sub.L, the photocurrent I.sub.x1 and I.sub.x2 outputted from the electrodes X.sub.1 and X.sub.2 are represented as follows: ##EQU1## and further, since the resistance between X.sub.1 and X.sub.2 is maintained at uniform distribution, the following respective equations are valid between the resistance between X.sub.1 and X.sub.2 and the length L: EQU R.sub.x1 +R.sub.x2 =L EQU R.sub.x1 =X EQU R.sub.x2 =L-X\nAccordingly, the signals taken out from the respective electrodes can be represented by L and x as follows: ##EQU2## Thus, the informations of the incident position of light and light intensity can be obtained at the electrodes of X.sub.1 and X.sub.2.\nFurther, by taking a ratio of the difference between I.sub.x1 and I.sub.x2 to the sum thereof and representing it as P, the following equation can be obtained: ##EQU3## and the positional signals irrelevant with light intensity change can be obtained continuously corresponding to x=0 to L as follows: ##EQU4##\nHaving described about the one-dimensional case, the two-dimensional case may also be considered similarly, and the positional signals can be determined from the block diagram of the actuation circuit as shown in FIG. 40.\nHere, from the actuation principle of PSD, when light is introduced at two or more points, the positional signal weighted in proportion to the respective light intensity is obtained. Also, in the case when the light flux is expanded, the positional signal corresponding to the gravitational center of light intensity can be obtained.\nOn the other hand, there has been known in the art a device for forming monomolecular built-up films in which monomelecular films are transferred onto a substrate by laminating them one by one according to the monomolecular layer film built-up method called the Langmuir-Blodgett method (hereinafter called LB method) named after the inventors (see Shin Jikken Kagaku Koza (New Course of Experimental Chemistry), Vol. 18, pp. 498-507, Maruzen).\nThe above device is constituted schematically by a liquid tank containing a liquid, a film-forming frame which is floated so as to divide the liquid surface into two and is capable of two-dimensional piston movement, a driving means for moving the film-forming frame, a surface-pressure-measuring instrument for measuring the surface pressure of the monomolecular film spread on the liquid surface and a substrate holder for moving the substrate held vertically relative to the liquid surface. Formation of a monomolecular film and transfer thereof onto a substrate are practiced as described below.\nFirst, with the film-forming frame being kept at one side of the liquid tank, a solution of a film-forming substance dissolved into a volatile solvent such as benzene, chloroform, etc. at a concentration of, for example, ca. 5.times.10.sup.-3 mol/liter is added in several drops onto the liquid surface by means of a fountain pen filler, etc. When the solution is spread over the liquid surface and the solvent is evaporated, the monomolecular film remains on the liquid surface.\nThe above monomolecular film exhibits the behaviour of a two-dimensional system on the liquid surface. When the surface density of molecular is small, it is called as a gas film of two-dimensional gas, and the state equation of two-dimensional ideal gas is valid between the occupied area per molecule and the surface pressure.\nSubsequently, by increasing the surface molecular density by narrowing the region of the liquid surface carrying the spread molecular film by moving gradually the film-forming frame, the interaction between molecules is intensified, whereby the gas film changes via a liquid film of two-dimensional liquid to a solid film of two-dimensional solid. When such a solid film is formed, the molecules become oriented fairly regularly and the film has highly ordered characteristic and uniform ultra-thin film characteristic. And, by moving vertically the substrate with the substrate holder, the monomolecular film which has become said solid film can be attached and transferred thereon. Also, by transferring the monomolecular film for plural times on the same substrate, a built-up monomolecular film can be obtained. As the substrate, there may be employed, for example, a glass, a synthetic resin, a ceramic, a metal, etc.\nIn order to practice the transfer operation under the preferable state of the monomolecular film for transfer onto the above substrate, the surface pressure of the monomolecular film is measured. The surface pressure preferable for transfer is generally accepted to be 15 to 30 dyn/cm. Outside this range, the alignment or orientation of the molecules may be disturbed, or otherwise peel-off of the film will readily occur. However, in a special case, for example, depending of the chemical structure of the film-forming substance, the temperature conditions, etc., preferable values of the surface pressure may be outside the above range, and therefore the above range may be a tentative measure.\nThe surface pressure of the above monomolecular film is measured automatically and continuously by means of a surface-pressure-measuring instrument. As the surface-pressure-measuring instrument, there is one in which the method for determining the difference in surface tension between the liquid surface covered with no monomolecular film and the liquid surface covered with monomolecular film is applied, or one in which the two-dimensional pressure applied on the film-forming frame which will float by partitioning the liquid surface into the liquid surface covered with no monomolecular film and the liquid surface covered with monomolecular film, etc., each having specific features. The occupied area per molecule constituting monomolecular film and the degree of the change are also generally measured together with the surface pressure. The occupied area and its change can be determined from the movement of the film-forming frame.\nThe movement is controlled based on the surface pressure of the monomolecular film measured by means of the above measuring instrument. More specifically, the driving means for moving the film-forming frame is controlled based on the surface pressure of monomolecular film measured by the surface pressure measuring instrument so that the monomolecular film can constantly maintain a predetermined surface pressure selected within the preferable range for transfer operation. The movement control of the film-forming frame is continuously made not only until initiation of the transfer operation of monomolecular film after dropwise addition of the film-forming substance, but also during the transfer operation continuously. For example, in the transfer operation, as the monomolecular film is transferred onto the substrate, the surface density of the monomolecular-film-forming molecules on the liquid surface will be lowered, whereby the surface density will be also lowered. Accordingly, by moving the film-forming frame, the spreading area of the monomolecular film is narrowed thereby to maintain the constant surface pressure by correction corresponding to the lowered surface pressure.\nHowever, when a PAS device, PTR device or PDS device itself is to be used for measurement of an extremely thin material under a specific environment such as a thin film of monomolecular film being spread on liquid surface, problems arise such as difficulty of measurement, liability of declining precision, sensitivity of measurement because of the location on a liquid surface of the extreme thinness of the measured material.\nPAS devices can be classified into the microphone system and the piezoelectric element system depending on the kind of the detector. In the case of the microphone system, a sample is required to be placed in a hermetically sealed sample chamber, while in the case of the piezoelectric element system, arrangement of the detector and the sample is restricted. Thus, either one of them is not suitable for measuring a thin film as spread on a liquid surface.\nIn the case of the collinear type PDS device, since the probe light passes through the measuring site of the thin film where the greatest refractive index change occurs by irradiation of exciting light, there is the advantage that a relatively greater positional change of probe light can be obtained on the detector. However, in the collinear type PDS device, since the probe light transmits through the thin film which the material to be measured, it receives at the same time the influences from both the liquid phase and the gas phase of which refractive index varies based on exciting light absorption of the thin film. Accordingly, no correct measurement is possible unless the variation in refractive index between the liquid phase and the gas phase is considered, whereby there is the problem that measurement of high precision becomes extremely difficult.\nAlso, when the material to be measured is a thin film as in the present invention, the refractive index change of the surrounding caused by absorption of exciting light is small. Accordingly, in the case of a transfer type PDS device, it is required that the probe light be brought close to the thin film on the liquid surface so that it can pass through the region where refractive index change as great as possible occurs. Particularly, since the detector of the PDS device has the property of detecting the position of the \"gravitational center\" of the light-receiving intensity, and it is preferred that the flux center of the probe light with strong light intensity be nearer to the thin film. However, it is not possible to access the center of the probe flux nearer than its radius. Therefore the probe flux center with strong light intensity is liable to pass through the region apart from the thin film, resulting in little refractive index change. This makes a high precision and high sensitivity measurement difficult.\nIn PDS measurement, the convergence state of the fluxes of exciting light and probe light have influences on sensitivity, and the relative positional relationship between exciting light and probe light will affect sensitivity greatly. Accordingly, it is necessary to monitor the irradiation positions of exciting light and probe light and their relative positions and, in prior art, scattered light from excitation light and probe light have been monitored by visual observation with a microscope, but such a method will give undesired external disturbance such as vibration of liquid surface. Also, in the case of a very thin or homogeneous film, the scattered light may be sometimes undetectable. Further, in carrying out measurement of light absorption characteristics, the wavelength of the light used is not limited to a visible light region and, for example, in the case of using IR-ray, there has been the drawback that the irradiation position cannot be confirmed by visual observation.\nFurther, concerning measurement of a material spread on the liquid surface such as monomolecular film, noise by wobbling of the liquid surface brings about lowering in measurement precision.\nOn the other hand, as described above, various fine controls are required for obtaining a monomolecular built-up film. However, many experiments must be conducted in order to search the optimum condition, and it can be confirmed only indirectly by the surface pressure, etc. whether the monomolecular film on the liquid surface is in the suitable condition for built-up or not, and thus precision is insufficient. This can be improved by grasping directly the physical properties of the monomolecular film on the liquid surface by means of PAS, PTR or PDS device, but due to the problems as mentioned above, such a demand cannot be met under the present situation."}
{"text": "1. Technical Field\nThe present invention relates in general to remote control of consumer electronic devices and in particular to computer-based programming of recording in remote controlled record/playback devices. Still more particularly, the present invention relates to utilizing a computer to control record programming in remote controlled record/playback devices in connection with a server push mechanism.\n2. Description of the Related Art\nControl over typical consumer electronics devices which utilize infrared remote control interfaces is typically extremely limited. In particular, programming such devices to record or play at specific times, etc. is generally constrained by the capabilities built into the device hardware. Programming such remote controlled devices may also be confusing, as has been frequently noted in connection with video cassette recorders (VCRs) and more recently in connection with unitary compact disc (CD) players, cassette recorders, and AM/FM receivers (often colloquially referred to as \"boom boxes\"). Flexibility over program control is not generally found in such devices.\nThe advent of set-top boxes, also referred to variously as a network computer (NC), Internet appliance, Web PC, Internet access device, browser box, etc. (all names for the same concept with minor variations), has introduced a new consideration. Frequently such set-top boxes are themselves capable of being controlled by an infrared remote control. Additionally, the purpose of the devices generally results in the devices being situated in close proximity to a television, VCR, or other remote controlled electronic device, often within the same entertainment center. Much of the same hardware and control software utilized to interface an infrared receiver to a set-top box could be reused to interface the set-top box to an infrared transmitter.\nIt would be desirable, therefore, to provide a method of controlling programmable electronics devices utilizing the infrared remote control interface under the control of a set-top box or other data processing system. It would further be advantageous to provide such control in a manner consistent with the network technique known as \"server push.\""}
{"text": "Field of the Invention\nThe present invention pertains to free weight lifting apparatus and securing freeweights to a barbell especially with respect to dumbbells."}
{"text": "Legacy system computer applications that have been executing on mainframe computers, in some cases, for decades, continue to play a vital role in American commerce, education, and industry. Such software, in fact, has vastly outlived the time period during which its creators believed it would be useful. Doubts about the truth of the foregoing proposition vanish in the face of the billions of dollars that have been (or will be) spent by U.S. companies and the United States government in correcting computer code that is not Year 2000 compliant. The vast bulk of this code is substantially more than five years old.\nAdded to this vast mass of legacy systems are modern systems executing in modern environments, such as the client/server environment. These systems, like the legacy systems, suffer from the same basic limitation--the business logic and associated data are locked up in architecture-specific systems. Those using different operating systems running on computers powered by entirely different central processing units cannot effectively access the data and business logic residing on these conventional systems.\nThese conventional systems (both the legacy systems and the more modern systems) are not only important, they also represent substantial investments by companies. They are the product of billions of dollars of programmer compensation and untold millions of hours of business planning, strategic design, and work flow description. In many respects, the code embodied in these conventional systems describes the operations of everyday commerce in this country.\nToday's technology culture differs widely from the environment that existed when most of these conventional systems were created. Access to computers is no longer limited to an elite handful of technicians with dedicated terminals hooked to gargantuan boxes located in specially cooled rooms. Today, virtually every worker has a personal computer sitting on his/her desk. Others, on the go, carry their computers in brief cases that connect to the home office via telephone lines. Consumers, from their dens and living rooms, are buying and selling stocks, checking their bank accounts, and (in some cases) telecommuting. The need to get the data and business logic out of the hands of the specially trained technicians and into the hands of every day people is great. The demand is the same, and the problem is the same, whether the important data and business logic resides on a legacy system or on a modern client/server system. Accessing and interacting with it from the outside, an external system, is a huge chore.\nThe demand is great, but, even more importantly, it is growing exponentially--with the explosion in popularity of the Internet and, especially, the World Wide Web. The number of trademark applications alone with terms like \"e-biz,\" \"e-business,\" \"e-commerce,\" and the like is staggering. Clearly, commerce has moved to the net. Virtually no television advertisement and certainly no television news program fails to include a reference to its Internet or World Wide Web address.\nCoupled with the popularity of the Internet has been another trend that has silently but markedly transformed the landscape of commerce and, indeed, everyday life in the United States. This is the trend toward self-service. Full service filling stations, for example, have almost been relegated to a dim memory. ATM machines have replaced tellers, and even some banks charge an extra transaction fee for using a human teller, as opposed to an ATM machine. The Internet itself is the ultimate expression of the trend toward self-service. Banking, book buying, car and home shopping, teaching, and even church services are available when the cyber-consumer wants it and without the presence of any human tellers, brokers, clerks, teachers, professors, registrars, or ministers.\nRegardless of the computer system on which the application is running, the same basic disability exists. The data and business logic are locked up in an architecture-specific format. Almost universally, raw data is stored in a format shaped primarily by storage constraints, e.g., in relational database tables, and by retrieval considerations, e.g., indexes. This data is presented in business-useful human-friendly form only when acted upon by the architecture-specific computer application that carry out instructions based on business logic. In other words, the data resides in one generally useless format and is put in useful form only when acted upon by a separate computer application, which application is typically architecture-specific in terms of its functionality. Thus, whether the information and business logic are locked up inside a venerable legacy system or are stored in client/server systems in database tables, systems on the outside have difficulty in reaching into these conventional systems for not just the raw data, but the data in a form and format that has been filtered, selected, organized, and processed by intelligence that embodies an organization's business logic.\nThe problem, as noted above, is not just converting the data from one system to another. Difficult as that problem is, converting the raw data from a format on one magnetic medium to a different format on another medium has been done before. As noted above, however, conversion of the data does not solve the problem. The business logic that sorts, parses, selects, combines, performs operations upon, and presents this data in useful form is left behind. The raw data, even if converted, is just sitting there. It is no longer part of the system. It is just plain data.\nIn order to make this data, even after it is converted into a form that can be processed by an external system, of any use, it has to be combined with the business logic, e.g., combined with computer code that sorts, parses, selects, combines, presents, and otherwise operates on the data in ways that are meaningful to the business or institution that owns the data. That business logic already exists in the executable code still residing on the conventional system. The business logic needs to be duplicated or replicated in a form that the external system can use.\nThe process of reinventing the business logic is inefficient, time-consuming, and (in many cases) ineffective. It is inherently inefficient to operate two parallel systems--the conventional system and the external system. The conventional system represents a huge investment of time and money--an investment likely to have painfully escalated with the cost of making these conventional systems Year 2000 compliant. Reinventing these programs in another environment--the environment of the external system--hardly makes sense now. Moreover, the programmers who best know the organization's business logic, as it is embodied in the programs executing on the conventional system, are probably not the people who are best equipped to recreate the business logic in the new (external) environment. Almost by definition, they are trained on and have developed expertise in the environment of the conventional system, not the external system. Thus, if the business logic is to be recreated on the external system, those with the best knowledge of the business logic will need to be retrained in the language of the external system or, in the alternative, the job of recreating the business logic in the external system will be handled by those with no experience in the organization's business logic. Both approaches fall far short of the ideal.\nAn alternative to the solution described above--translation of the data and recreation of the business logic--would be to graft external system awareness into the conventional system. This approach is best illustrated by the efforts of some to transplant HTML-aware routines, libraries, and tools into legacy systems. Modern flavors of RPG, COBOL, Fortran, and others sport new web awareness tools and extensions. This approach has both advantages and disadvantages.\nThe main advantage of this approach is that it allows the legacy system programmer to continue to function in the environment in which he/she feels comfortable. These legacy system programmers do not have to become web gurus or become proficient in Java, HTML, and other web-based or web-aware languages.\nThe disadvantage involves the same problem of duplication mentioned above. The legacy system programmers who have just spent the last two or three years plowing through lines of codes that haven't seen daylight in decades to fix the Y2K problem now have to go through all those lines of code all over again, rewriting them to become not Y2K compliant, but web compliant. In other words, the systems will all have to be modified and/or rewritten using these new web-aware tools.\nThe applicant's invention permits organizations, companies, and institutions with legacy and other conventional systems to make not only their raw data but their business logic available to external systems without having to create a duplicate system in a parallel external universe and without requiring major code revisions in the conventional systems themselves. The applicant's novel approach takes the data after business logic has been applied within a running program and channels it into the external system environment intact, where the intact data is reconfigured into data that is understandable by the external system.\nMoreover, this packaging and processing of intelligent data (data acted upon by the conventional system's business logic) takes place at the \"code level.\" To understand what is meant by the phrase \"code level,\" it is important to understand the preexisting \"translation\" technology. Because of the inherent differences between legacy system data formatting and language conventions and those popularized by the PC revolution, the most popular translation form is commonly referred to as \"screen scraping.\" This moniker is actually very descriptive of the process that is in fact used. The data and business logic inherent in a legacy system are presented in meaningful, humanly comprehensible form when it is printed on the computer screen. It does not really matter whether the computer screen is hooked to a mainframe computer or a PC connected to the World Wide Web. What meets the human eye is equally comprehensible. The screen scraping technique, therefore, basically uses a series of algorithms to read and/or intercept the display function, scraping, as it were, that human-readable information and piping the result into an image projected on the screen of an external system.\nThe screen scraping technique has numerous problems. First, it is very computer intensive and very, very slow. Second, it is, of necessity, screen dependent. Screen scraping requires one to model the flow of a particular application. Modeling program flow entails a huge number of complexities, particularly where user input is involved or error recovery is required. When one is literally scraping the screen, it is necessary to anticipate which screen is coming next in program flow. Third, with screen scraping, it is assumed that the screens will be static. If, however, circumstances require modifications to any of the screens, regardless of the reason, further modification to the screen scraping system will be required to avoid errors. Fourth, screen scraping is extremely difficult when the business logic requires one to gather pieces of data from different screens and assemble all those bits and pieces of information together. Fifth, screen scraping is not scalable at all. It is inflexibly dictated by the architecture of the conventional system and the external system. Sixth, the result of screen scraping in the web environment is not particularly desirable. Without substantial further processing, its presentation is less than fully satisfying to the Web-experienced user. Indeed, without additional processing, it may be incapable of interacting with the user. Although, for discussion purposes, the foregoing has focused upon screen-scraping of legacy systems to make them Web-aware, the same issues exist regardless of the types of conventional systems and external systems involved.\nRather than scrape the screen, Applicant's invention involves taking a snapshot of the in-memory data, as it is being processed by the software applications residing on the conventional system. By taking a snapshot, not of the screen, but of the in-memory, real time image of the data, Applicant's invention operates at the code level and, hence, is highly efficient. It is also highly scalable. To interact with the business logic and data residing on the conventional system, conventional system programmers can continue to call upon their conventional system tools and expertise.\nThe Applicant's invention has none of the limitations that have plagued the prior art solutions. Indeed the disadvantages of the prior art offerings are matched item-by-item by strengths of the applicant's invention. Some of the advantages of the Applicant's invention are that it allows proven functionality to continue, using mature code that has been thoroughly debugged and is reliable. New applications do not have to be written in order to make the business logic and precious data available for interaction by persons operating systems different from the conventional systems. The Applicant's invention is not screen-dependant. Rather, the execution-time in-memory representation of the data, as acted upon by the business logic, is translated and made available to the external systems. At the same time, information from an external system, may be fed into and processed by the conventional system by means of Applicant's invention. In addition, Applicant's invention is scalable to an almost unlimited degree in that, by making the data and business logic available in any external system, including markup language systems, such data and business logic may be accessed and used on any platform. The data and business logic become wholly and completely platform independent. Another benefit is in the nature of human cost. To satisfy the demand for external system access to conventional system data and business logic, including markup language versions of the data and business logic now trapped inside conventional systems, it would be necessary to rewrite the application, with the attendant human costs in terms of both programming and training or retraining. The Applicant's invention also has the advantage of making conventional system data and business logic available to external systems without making huge demands on processor and computing resources."}
{"text": "1. Field of the Invention\nThe present invention relates to a parallel operational processing device, and particularly to a construction of a parallel operational processing device having a semiconductor memory and an ALU (arithmetic and logic unit) integrated together.\n2. Description of the Background Art\nIn a field of portable equipments and in application of image processing, it has been recently demanded, due to improved performances, to process a large amount of audio or video data at a higher speed. For processing a large amount of data, a dedicated DSP (Digital Signal Processor) is generally used in many cases. When the dedicated DSP is used, hardware thereof is fixedly set, and a data bit width and processing contents are restricted. For flexibly accommodating for various applications, it is now required to execute such processing with a programmable processor that can change its functions through software. Particularly, in view of reduction in size, it is important in such a data processing system to achieve fast processing with a small area and low power dissipation.\nIn the processing of audio and image data, a plurality of sets of data are frequently subject to the same arithmetic or logic operation. For such a processing application, therefore, SIMD (Single Instruction Stream Multiple Data Stream) processor is often used. In such SIMD processor, processor elements each constructed by a register file and an ALU are arranged in parallel, and are operated in parallel by the same instruction. These processor elements in parallel are 256 in number, enabling parallel operational processing on 256 sets of data, as disclosed in Reference 1 (Japanese Patent Laying-Open No. 2003-186854).\nIn the SIMD processor, each processor element stores multi-bit data in a register file. The ALU receives corresponding multi-bit data from a register in the file, and executes operational processing. The result of the operational processing is stored back into a register of the corresponding register file. Therefore, the bit width of the operation data to be processed depends on a bit width of the ALU or register. As disclosed in Reference 1, the change in bit width of the operation data can be handled by changing the number of the registers used in the operation in the register file. However, when 16-bit operation data is to be stored using, e.g., an 8-bit register, it is necessary to set an upper-bit register and a low-bit register, which in turn are successively switched for performing an operational processing. Therefore, when multiplication processing is performed according to a Booth's algorithm by obtaining a partial product and then adding the partial products to obtain a final product, such a problem arises that switching between the registers is extremely complicated for aligning digit positions of the operation data.\nAlso, the processing procedure must be changed depending on switching and non-switching of the registers. Accordingly, the processing procedure must be changed according to the bit width of the operation data, and it is impossible to deal with the change in bit width with flexibility.\nThe SIMD processor performs the parallel arithmetic or logic operation on multi-bit data, and the ALU is a multi-bit ALU and therefore requires a large layout area, which impedes reduction in size.\nFor overcoming such problem of the SIMD processor as described above, the group of the inventors has already proposed a parallel operational processing device achieving a fast arithmetic/logic operation (e.g., by Japanese Patent Application No. 2004-171658 and Japanese Patent Application No. 2005-143109 (U.S. Ser. No. 11/148,369)), the contents of which are incorporated herein by reference. The parallel operational processing device has a basic construction in which a memory cell array is divided into a plurality of entries, and one-bit ALU is arranged corresponding to each entry. Each entry stores operation data. Each ALU performs arithmetic or logic operation in bit serial manner. For example, data bits are read from corresponding entries in memory cell arrays that are placed on the opposite sides of the ALUs, and are transferred to the corresponding ALU, which in turn execute the operation on the received data bits and stores the result of operation in a predetermined entry.\nThe arithmetic/logic operation is performed on multiple words in bit serial manner (word parallel and bit serial manner). Therefore, an operation on each data item requires much time, but fast processing can be achieved by increasing the number of entries and increasing the degree of parallelism. For example, in an operation environment requiring one machine cycle for each of read, write and operational processings, one-bit arithmetic/logic operation requires three machine cycles. Therefore, the processing of 8-bit data requires twenty-four (=3·8) machine cycles. However, when the entries are provided, e.g., 1024 in number, operation on 1024 sets of data can be completed in twenty-four machine cycles. The operation data is generally 32 or 64 bits in bit width. Therefore, fast operations can be achieved by increasing the number of entries.\nSince the arithmetic/logic operation is performed in a bit serial manner, the change in bit width of the operation data can be accommodated for by changing an accessing address range of the entry. The change in processing procedure such as register switching is not required, and the change in bit width of the operation data can be easily made.\nThe ALU is a one-bit ALU, and can have a small layout area. When bit lines of a memory cell array are used as entries, the ALUs can be arranged corresponding to the entries, respectively, and it is possible to achieve the parallel operational processing device having a small layout area and an extremely high operation parallelism.\nFor further increasing the operation parallelism, it is necessary to arrange more efficiently the ALUs and peripheral circuitry of the memory cell array. However, a region where the ALUs and the peripheral circuits are determined depends on a pitch of memory cells. Since a layout rule for the memory cell array and the pitch of memory cells are set in advance, further improvement is required for efficiently arranging the ALUs and the peripheral circuitry with a layout area reduced further.\nThe memory cell array is formed using, e.g., an SRAM (Static Random Access Memory) not requiring refreshing. Since the arithmetic/logic operation is performed on data in bit serial manner, when read modify operation is executed in which the reading of data and the writing of an operation result are performed in one cycle, in order to speed up the operational processing, the operation frequency of the SRAM that performs reading and writing of data in one machine cycle becomes higher than that of the ALU. Consequently, the operation frequency of the SRAM limits the operation speed of the parallel operational processing device, and therefore the operating manner of the SRAM must be further considered for achieving faster arithmetic/logic operations.\nFurther, a higher operation frequency increases power consumption. For increasing the operation speed, it is necessary to increase layout area of the elements for increasing a current driving capability, resulting in an increased power consumption. Further, the number of entries must be increased for increasing the parallelism degree of operations. In such case, the circuits operating in parallel increase in number, and the current consumption increases. Therefore, further devising is required also in view of the above for increasing the degree of operation parallelism and the operation speed without increasing the power consumption."}
{"text": "User authentication is implemented to ensure that a user is who they purport to be. For example, a user of a computing device can attempt to log into a back-end system. User authentication can be implemented to ensure that the user is authorized to log into the back-end system. In some examples, authentication can be achieved based on user credentials (e.g., username password).\nBiometric authentication enables users to be authenticated based on physical features that are unique to the user (e.g., facial features, voice). Biometric authentication is advantageous in that it is not easily spoofed, and that it does not require the user to remember credentials (e.g., username, password). Forms of biometric authentication, however, might not be appropriate for a particular environment that the user is in. For example, facial recognition and/or voice recognition might not be appropriate forms of biometric authentication, while the user is in a movie theater."}
{"text": "The present invention relates to a device for closing and coupling luggage and the like, particularly for motorcycles.\nRecently it has become popular to apply on the rear rack of motorcycles, both low- and high-powered ones, luggage or rigid bags which are applied by means of bracket-like elements provided on the case to engage coupling pins provided on the racks.\nSuch luggage are disclosed in prior patents in the name of the same Applicant, such as for example Italian patents 1,154,627 and no. 201,809, which are assumed included herein as reference.\nThe greatest problems of current motorcycle luggage, such as carryall cases, bags etc. are linked to two factors; namely, first of all the need to prevent theft of the cases or rigid bags, and secondly the need to provide easy opening of the case, combined with easy uncoupling of the case from the motorcycle, obviously by the user who owns a related opening key.\nFurthermore, due to reasons which are also linked to the streamline of the motorcycle, it is important to eliminate the fixed external handle of the case, which currently protrudes from the body of most carryall cases or rigid bags installed on motorcycles."}
{"text": "1. Field of the Invention\nThe present invention relates to a data processing apparatus, data processing system, and access area control method and is particularly suitable to a data processing apparatus in which devices having different numbers of address signal lines are connected via an address bus.\n2. Description of the Related Art\nConventionally, an information processing system in which a plurality of devices such as central processing units can access (can read out data from and write data in a memory) a memory area (memory space) composed of one or a plurality of memories exists. In one information processing system of this type, a plurality of devices are not connected to a memory bus for accessing a memory, but, as shown in FIG. 6, only a device 61 is connected to a memory bus 67. In this conventional information processing system shown in FIG. 6, another external device (not connected to the memory bus 67) 63 accesses a memory 62 via the device 61 connected to the memory bus 67.\nAs shown in FIG. 6, this information processing system comprises a central processing unit (CPU) 61, SDRAM (Synchronous Dynamic Random Access Memory) 62, external bus master 63, ROM 64, RAM 65, and input/output buffer (I/O) 66.\nThe CPU 61 has a CPU core 70 for performing arithmetic processing, an SDRAM controller 71 for performing data write and read to the SDRAM 62, and an external bus controller 72 for exchanging data with the external bus master 63 and the like. The CPU core 70, SDRAM controller 71, and external bus controller 72 are connected to an internal bus 73.\nAlso, the SDRAM 62 and the SDRAM controller 71 are connected via an SDRAM bus 67, and the external bus master 63 and the external bus controller 72 are connected via a general-purpose bus 68. This general-purpose bus 68 is connected to the ROM 64, RAM 65, and input/output buffer 66.\nReferring to FIG. 6, examples of signals exchanged between the CPU 61 (external bus controller 72) and external bus master 63 through the general-purpose bus 68 are indicated by the dotted lines. These signals includes, e.g., a bus request signal BREQ#, a bus grant signal BGRNT#, an address signal A<31:0>, a data signal D<63:0>, a SDRAM select signal SDSEL#, and a read/write signal RW#.\n“(Signal name)#” indicates that the signal is activated at low level, and “(signal name)<X:Y>” indicates that the signal is transmitted through a plurality of signal lines in parallel. For example, the address signal A<31:0> indicates that this address signal is transmitted through 32 address signal lines in parallel with a width of 32 bits. The same rules apply to “(signal name)#” and “(signal name)<X:Y>” hereinafter.\nIn the information processing system shown in FIG. 6, operations when the CPU 61 (CPU core 70) and external bus master 63 access the SDRAM 62 will be explained below. In the following explanation, each address bus width of the general-purpose bus 68 and internal bus 73 is 32 bits.\n<When CPU 61 (CPU Core 70) Accesses SDRAM 62>\nThe CPU core 70 outputs the address signal (address value) and read/write signal to the internal bus 73. The SDRAM controller 71 detects that the address signal supplied via the internal bus 73 indicates a memory space in the SDRAM 62. On the basis of the supplied address signal and read/write signal, the SDRAM controller 71 controls the SDRAM bus 67 to access the SDRAM 62. This allows the CPU core 70 to access the SDRAM 62.\n<When External Bus Master 63 Accesses SDRAM 62>\nFirst, the external bus master 63 activates the bus request signal BREQ# and requests the CPU 61 (external bus controller 72) to issue the bus use right (to be simply referred to as the “bus right” hereinafter) of the general-purpose bus 68. In response to this request, the external bus controller 72 activates the bus grant signal BGRNT# and gives the external bus master 63 the bus right.\nWhen acquiring the bus right of the general-purpose bus 68, the external bus master 63 activates the SDRAM select signal SDSEL# which indicates access to the SDRAM 62. In addition, the external bus master 63 outputs the address signal A<31:0> and read/write signal RW# to the general-purpose bus 68.\nOn the basis of the SDRAM select signal SDSEL#, the external bus controller 72 detects the access request from the external bus master 63 to the SDRAM 62. The external bus controller 72 outputs, to the internal bus 73, an address signal and read/write signal corresponding to the supplied address signal A<31:0> and read/write signal RW#, respectively. In addition, on the basis of these address signal and read/write signal supplied via the internal bus 73, the SDRAM controller 71 controls the SDRAM bus 67 to access the SDRAM 62.\nAs described above, access from the external bus master 63 to the SDRAM 62 is realized via the CPU 61. That is, access from the external bus master 63 to the SDRAM 62 is executed by a “transparent mode access function” by which the CPU 61 accesses the SDRAM 62 via the SDRAM bus 67 in response to a request from the external bus master 63.\nIn this “transparent mode access function”, in response to a request from an external bus master (the external bus master 63) connected to an external bus (the general-purpose bus 68), a first bus master (the CPU 61) accesses a resource (in the above explanation, the SDRAM 62) which the external bus master cannot directly access and which is mapped in the memory space of the first bus master, thereby executing access from the external bus master to the resource. The access from the external bus master to the resource mapped in the memory space of the first bus master, which is realized by the transparent mode access function described above, is called “transparent mode access”.\nIn the above conventional information processing system shown in FIG. 6, however, if the address bus width of the external bus master 63 with respect to the general-purpose bus 68 is smaller than the address bus width of the CPU 61, transmission mode access can be performed only to a partial area fixed in accordance with the address width of the external bus master 63 and cannot be performed for an arbitrary area. For example, when the address bus width of the external bus master 63 is 8 bits, the external bus master 63 can perform transmission mode access only to an area from “0000 0000 h” to “0000 00FF h” in the memory space of the CPU 61 (“h” indicates the hexadecimal notation, and the same shall apply hereinafter). Therefore, to allow the external bus master 63 to perform transparent mode access to an arbitrary area in the memory space, a new external circuit (device) must be formed. This increases the external circuit scale and complicates the wiring between the devices.\nAlso, if the address bus width of the external bus master 63 with respect to the general-purpose bus 68 is equal to the address bus width of the CPU 61, the whole memory space of the CPU 61 can be subjected to transparent mode access. Accordingly, if an access inhibited area which inhibits access from the outside is necessary in the memory space, an external circuit (device) for limiting input address signals must be formed. This increases the external circuit scale and complicates the wiring between the devices."}
{"text": "This invention relates to a new propulsion system for satellites. Specifically this invention relates to a reduced toxicity satellite fuel that can be used for both the maneuvering and station-keeping propulsion systems of a satellite.\nCurrent satellite propulsion systems typically use nitrogen tetroxide with hydrazine in bipropellant class thrusters for maneuvering propulsion and use hydrazine in monopropellant class thrusters for stationkeeping propulsion. Unfortunately these satellite propellants are highly toxic and therefore, require special handling, transportation, and storage mechanisms, which add substantial cost to the deployment of satellites.\nOne of the goals of NASA\"\"s Discovery Program for new planetary exploration missions, is to substantially reduce total mission cost while improving performance. The performance and cost of the on-board propulsion system for satellites can be a significant factor in obtaining the highest possible science value per unit cost.\nConsequently there exists a need for lower cost reduced toxicity fuels with thrust per unit mass flow and density characteristics that are sufficient to replace prior art toxic fuels. Reduced toxicity fuels have not been used in the past, due to the fact that candidate fuels are not hypergolic. In other words, liquid reduced toxicity fuels will not spontaneously react with an oxidizer to begin the combustion process as in prior art fuels such as hydrazine.\nThus, to produce a bipropellant satellite thruster for use with a reduced toxicity fuel, there further exists a need for the thruster to have an ignition element consisting of decomposing elements for decomposing a reduced toxicity propellant into hot gases. These hot gases, like hypergolic toxic liquid fuels will spontaneously react with an oxidizer and begin the combustion process.\nIn addition to being used with bipropellant class thrusters, there is a further need for this reduced toxicity fuel to be used with monopropellant class thrusters. As a monopropellant, the reduced toxicity fuel must have a molecular structure that will decompose into low molecular weight gases without the formation of a solid constituent such as graphite. These monopropellant thrusters must also contain decomposing elements for reforming the reduced toxicity fuel into propellant gases. Satellite fuels that can be used as both a monopropellant and a bipropellant are referred to as dual-mode fuels.\nIt is an object of the present invention to provide a reduced toxicity propellant for use in satellite propulsion.\nIt is a further object of the present invention to provide a satellite thruster with the ability to catalytically decompose a reduced toxicity propellant into hot gases.\nIt is a further object of the present invention to provide a satellite thruster with the ability to decompose a reduced toxicity propellant into hot gases with a fuel cell reformer.\nIt is a further object of the present invention to provide a satellite thruster with a low weight plasmatron capable of decomposing a reduced toxicity propellant into hot gases without overheating and eroding portions of the plasmatron.\nIt is a further object of the present invention to provide a reduced toxicity dual-mode propellant that can be used in both bipropellant and monopropellant satellite propulsion systems.\nFurther objects of the present invention will be made apparent in the following Best Modes for Carrying Out Invention and the appended claims.\nThe foregoing objects are accomplished in one preferred embodiment of the invention by replacing the toxic fuel used in prior art satellite propulsion systems with a reduced toxicity liquid fuel such as methylamine. The thrusters in the present invention include a decomposing element for converting the reduced toxicity fuel into hot gases. These decomposing elements are included in both the monopropellant altitude control system (ACS) thrusters for stationkeeping and the bipropellant axial thrusters for maneuvering the satellite.\nIn the ACS thrusters, these decomposing elements are operative to decompose the reduced toxicity liquid propellant into propellant gases. In the axial thrusters the decomposing elements are operative to decompose the liquid reduced toxicity propellant into hot gases which auto-ignite with the second propellant in the combustion chamber of the axial thruster and thereby produce thrust when ejected through a nozzle. The difference between the thrusters is primarily their thrust class or the force generated during firing. The monopropellant ACS thrusters are in a smaller thrust class than the bipropellant axial thrusters because they are required to satisfy a minimum impulse-bit (thrust times time) requirement for precision pointing of the satellite.\nThe prior art uses a toxic propellant such as hydrazine in both the monopropellant ACS thrusters and bipropellant axial thrusters. Hydrazine is a hypergolic fuel, which means it will spontaneously react with an oxidizer such as nitrogen tetroxide in the liquid state thereby triggering the combustion process in prior art axial thrusters. Unfortunately, as discussed above, reduced toxicity propellants suitable for use with satellite propulsion are not hypergolic. Before the reduced toxicity propellants of the present embodiment will react with a second propellant, they must be decomposed into hot gases. These hot gases will auto-ignite with the second propellant and thereby begin the combustion process.\nPropellants can be decomposed by a number of different technologies, including the use of catalytic decomposing elements, fuel cell reformers, and plasmatrons. Each of these decomposing elements is suitable for different reduced toxicity propellants. For example, the amine, methylamine, the nitroparaffin, nitromethane, and the ether, ethylene oxide, can be catalytically decomposed. Alcohols such as methanol and ethanol, and saturated hydrocarbons such as methane can be decomposed with fuel cell reformers. Saturated hydrocarbons such as pentane and octane and jet engine fuels such as kerosene and JP-10 can be decomposed with a plasmatron. Other embodiments use unsaturated hydrocarbons such as 1-pentene, ring compounds such as cyclopropane, and strained ring compounds such as quadricyclane.\nIn the preferred embodiment of the invention the second propellant is an oxidizer such as nitrogen tetroxide, liquid oxygen, hydrogen peroxide, or oxygen difluoride. Although oxygen difluoride is highly toxic and must be handled as a mild cryogen on the ground, it represents a high performance option. Although hydrogen peroxide has a rather high toxicity, it has unique characteristics in that it is an unstable molecule that can be catalytically decomposed into hot oxygen rich gas. Thus hydrogen peroxide is suitable in use as both a monopropellant in the ACS thrusters and as an oxidizer in the axial thrusters.\nIn the preferred embodiment of the present invention the decomposing element of a thruster is always active decomposing the reduced toxicity fuel into hot gases. However, in alternate embodiments the decomposing elements could be used in an axial thruster to initiate the combustion process. Thereafter both propellants can be added directly to the combustion chamber and the decomposing element can be deactivated."}
{"text": "The present invention relates to an information display device for a camera, and in particular, to a liquid crystal display that displays exposure-related information and information that is imprinted into a film.\nMost cameras of present age are electronically controlled and comprise microcomputers as vital elements as well as peripheral circuits and peripheral units, wherein various functions are provided for easy exposing operation and special effects.\nFIG. 4 illustrates a circuit constitution of an electronically controlled camera according to the invention. To a microcomputer 11, ON/OFF signals are fed as control commands from various operation/control switches. These operation/control switches include contact switches 12 and 13 that are synchronous with an operation of a release button, and contact switches 14 and 15 that are synchronous with an operation of a zoom switch. Moreover, button switches 17 through 22 are provided to set an exposing mode, a display mode of a liquid crystal display 16, and a superimposing mode. Based on the ON/OFF signals from these operation/control switches, the microcomputer 11 performs various types of controlling. The peripheral circuits and units include a photometric IC device 23 and a photodiode 24 both for external luminance measurement; a shutter driving motor-control IC device 25 and a motor 26 both for driving and controlling the shutter; a focus lens driving control IC device 27 and a motor 28 both for shifting the focus lens to a focal position; a strobe control circuit 29 and a discharge tube 30 both for controlling flashing and discharging of the strobe; a focusing control circuit 31, an infrared light emitting diode 32, and a semiconductor-position detector (PSD) 33, each for detecting the distance to a subject; a zoom-lens driving control circuit 34 and a motor 35 both for shifting a group of movable lenses in the zoom lens to a position for an optimum focal distance; a film feed control circuit 36 and a motor 37 both for winding up or back a film; and a power circuit comprising a battery 38 and a power switch for a peripheral circuits.\nThese peripheral circuits respond to commands or control signals from the microcomputer 11, and return measurement signals and status signals to the microcomputer 11, whereby the microcomputer 11 performs programed controllings in order to acomplish the required exposing conditions. For example, once the release button is pressed, the microcomputer upon reception of an ON signal from the contact switch 12 controllingly performs photometry and range-finding, and then adjusts the position of the focus lens based on the resultant signals; and at the same time, the microcomputer verifies the shutter speed and chooses a status of the strobe. Upon reception of an ON signal generated by the actuation of a second-stage contact switch 13 on the release button, the microcomputer performs exposing by means of actuation of the shutter and controlled flashing of the strobe; and then, the film is wound up, and indications for film-winding status and film counter are updated at same time.\nPrior to picture-taking, by the manipulation of a drive mode switch 21, the present exposure mode rotates from regular exposing mode to single exposure mode, continuous exposure mode, and self start mode; and it is indicated on the display. Depending on the manipulation of a function mode switch 19, one of the multiple exposure mode, bulb mode, and interval mode can be selected. After a relevant mode is selected, selection of a data change position with a select switch 18 and various data settings, such as a number of times for multiple exposure and interval time with an adjust switch 17, are performed.\nNext, upon initiation of exposure control, appropriation of the power level of a battery 38 is to be checked. This checking procedure is achieved by A/D conversion of the voltage of the power switch 39, with the microcomputer 11 comparing with the reference voltage data. As described above, in a camera having battery checking means, check on the power level of the battery has been conventionally performed by comparing the power level with a specific reference voltage. Accordingly, even if the power level drops to a lowest permissible voltage, exposure control allows normal exposing. However, if the exposing mode is interval, a picture-taking operation possibly becomes impossible in the course of exposing. More specifically, in the interval mode, a predetermined number of frames are unconditionally subjected to exposing and film forwarding at predetermined, intervals. Therefore, even being at a proper level during first frame exposing, the power of the battery may drop to an impermissible level before exposing for the final frame. This voltage drop possibly results in immobilization or expersing defects such as being underexposed, and out of focus.\nOne possible measure to solve this problem is to set higher the lowest permissiblve voltage level. However, this measure judges that a battery comes to the end of its service life while it is still useful, and causes a waste of a useful battery.\nThe first object of the invention is to provide a battery check device that eliminates the waste of a useful battery and allows a reliable exposing in the interval mode.\nTurning the interval exposure control mode, exposing and frame forwarding are automatically performed at predetermined time intervals until a predetermined number of pictures are taken. For focusing control in the interval exposure, the focal length is either fixed to the initial position, or readjusted for each frame. The fixed focusing has a problem; pictures are readily out of focus when they are taken of a movable subject such as an animal. This problem can be solved by the readjusting focus control. However, with the latter focus control, if a subject is within a minimum exposure range, a warning is issued and exposure is inhibited as well as a regular exposing. Accordingly, in the interval exposing, if a subject is within the minimum exposure range, the corresponding frame is not exposed (one frame skipped); therefore, problems such as the cause of unexposed frame is not detected and an unattended interval exposing, occurs.\nOne possible measure to solve this problem is to perform exposing with focus unlocked. However, this measure poses another problem that a photographer may fail to notice that the camera is set at an out-of-focus position and performs an interval exposing. As a result, all or almost of all the frames can be exposed out of focus.\nThe second object of the invention is to provide an interval exposure control device that eliminates skipped frames while decreasing the possibility of out-of-focus exposure.\nExposure related information such as selections of different kinds of mode, data settings, film frame position, bulb time, and interval time are indicated on the liquid crystal display (LCD) 16. The LCD 16 also indicates data such as a date being imprinted on a film. Namely, in the vicinity of film surface at the inside of camera is provided an internal liquid crystal indication part (not shown in the drawings) for indicating information to be exposed on the film. The indication controlling system of the microcomputer 11 is so constituted that the internal liquid crystal indication part is adapted to indicate the information identical with that of the liquid crystal display 16. FIG. 5 illustrates a display pattern on the LCD 16 that comprises display segments for single exposure mode \"S\", continuous exposure mode \"M\", multiple exposure mode \"ME\", bulb mode \"Bulb\", interval mode \"INTV\", and that also comprises a numerical display section N for displaying present film frame position and multiple exposure timing, and a multi-digit numerical display TN for time display of interval timer as well as for date display.\nInformation display on the LCD 16 is controlled by the microcomputer 11, wherein a same numerical indicator selectively displays the date from the imprinting information or left time at the bulb mode; or elapsed interval time. In the case of a display pad which has a part commonly used for plural indications, the indication of information such as date and time, is switched based on the present status of exposure control.\nAs described above, in a prior art information display device, which selectively displays exposure related information and imprinting information, the display normally displays exposure related information, and switches to imprinting information to be also displayed on the internal liquid crystal indication part for exposing it on the film surface once picture-taking control starts (a contact switch 12 on the release button is turned ON). Accordingly, a photographer cannot monitor the exposure related information while picture-taking control is working. For example, in the case of picture-taking in the interval mode, the numerical display section of the LCD 16 switches to the imprinting date after picture-taking control starts, and a photographer fails to monitor the interval time elapsed.\nTo solve the above-mentioned problem, one possible measure is to provide a different display for each item of information instead of a dual-purpose display. However, this constitution inevitably increases the number of display segments to be controlled, and a larger number of display drivers; it creates a complicated and expensive display system. Another possible measure is to provide display modes being changed by a switch operation after the initiation of exposure control. However, this measure inevitably increases the number of switches required or requires more complicated operation by a photographer.\nThe third object of the invention is to provide an information display device that eliminates the problem in regards with switching indication of exposure related information and imprinting information."}
{"text": "1. Field of the Invention\nThe present invention relates to a microphone unit that converts a sound pressure (generated such as by voice) into electrical signals and outputs the resulting signals.\n2. Description of Related Art\nConventionally, microphone units are applied to voice input devices, examples of which are sound communication devices such as mobile telephones and transceivers, information processing systems, such as voice authentication systems, that utilize a technology for analyzing input voice and recording devices. Conventionally, as these microphone units, microphone units having various configurations are developed (for example, see patent document 1).\nFIG. 6 is a schematic cross-sectional view showing an example of the configuration of a conventional microphone unit. Among conventional microphone units, there is a microphone unit that includes, as with the microphone unit 100 shown in FIG. 6, a board 101 that has an opening portion 101a which is used as a sound passage, a MEMS (micro electro mechanical system) chip 102 that is mounted on the board 101 so as to cover the opening portion 101a, an IC (integrated circuit) 103 that is mounted on the board 101 and is connected electrically through a wire 105 to the MEMS chip 102, and a shield cover 104 that covers the MEMS chip 102 and the IC 103 mounted on the board 101.\nThe microphone unit 100 configured as described above converts, with the MEMS chip 102, sound signals transmitted through the opening portion 101a into electrical signals, amplifies the resulting electrical signals with the IC 103 and outputs the amplified signals through an unillustrated electrical pad to the outside.\nPatent document 1: JP-A-2008-92183"}
{"text": "1. Field of the Invention\nThe present invention relates to a body structure for vehicle bodies, and in particular to a front body structure forming a front portion of a passenger compartment for automobiles that are one of typical vehicles.\n2. Description of Related Art\nRecently various kinds of so-called hydro-forming have been proposed as one type of methods of forming metal members in the automobile industries. Hydro-forming is a molding or forming method to form a high tensile strength steel pipe to a given shape. The pipe member is placed in a die cavity and filled with a fluid such as water or oil. The pipe member is expanded by pressurizing the fluid until the pipe member is deformed to an external shape reflecting the contour of the die cavity. Such hydro-forming is known from, for example, Japanese a Unexamined Patent Publications Nos. 10-218017 and 11-235963 as applied to forming a reinforcement having a rectangular closed cross-section for a side sill or a center pillar of a vehicle body. In the hydro-forming, a pipe member for a reinforcement of a side sill is placed between upper and lower dies between which a cavity having a contour identical to an external shape of the reinforcement is formed, and then, pressurizing fluid is introduced into the pipe member so as to expand the pipe member until the pipe member is deformed to a shape reflecting the given contour of the cavity. The reinforcement thus formed is put in a follow side frame of a vehicle body and welding to the side frame. According to another example, a reinforcement having a rectangular closed cross-section, which is used to reinforce a center pillar extending downward from a roof side rail, is formed from a pipe member by the hydro-forming. The reinforcement is put in an internal space of the center pillar comprising a center pillar outer and a center pillar outer.\nThe hydro-forming can provides a pipe member with any external shape according to the given contour of the die cavity. By using this technology it becomes possible to obtain more lightened but more rigid structural members having various desired shapes as compared with similar members that are made by welding a plurality of stamped steel plates.\nGenerally a cowl structure located in front of a cabin or passenger compartment of an automobile that is a typical one of vehicles is an important structural part for safety because the cowl structure is not only contributory to ensure a body rigidity but also works as a base which mounts thereon a steering shaft and an air-bag that require high-accuracy and delicate adjustment in fitting position. However, the cowl structure needs a large number of parts that are usually complicated in shape, so that many restraints are put on designing and assembling the parts. Although, it is desirable to form the cowl structure from a reduced number of parts formed by the hydro-forming in light of simplifying assembling the parts to a cowl structure, nevertheless, a more reliable structure is needed because the cowl structure is of very importance in the terms of sufficient rigidity and safety of the vehicle body.\nAn object of the present invention is to provide a front body structure for vehicles which provides the front body with light weight and high rigidity.\nThe above object of the present invention is accomplished by a front body structure that comprises a dash panel separating a cabin or passenger compartment from the remaining part of a vehicle body, a dash panel member extending in a transverse direction of the vehicle body which is welded to the dash panel or another member disposed in the close vicinity of the dash panel, and an instrument panel member extending in the transverse direction in the close vicinity of the dash panel and supported at opposite ends thereof by side panels. The dash panel member and/or the instrument panel member are formed from a hollow pipe formed by hydro-forming. Each of the dash panel member and the instrument panel member has opposite ends welded to the side panels of the vehicle body, respectively.\nAccording to the preferred embodiment of the present invention, a pipe member as the instrument panel member has opposite ends welded to the side panels having door openings, respectively, and a generally U-shaped central portion opening upward and welded to a floor panel at a bottom thereof.\nAccording to another preferred embodiment of the present invention, a pipe member formed as the dash panel member has opposite end portions, each of which extends downward on the inner wall of the side panel along the door opening so that the end of the dash panel member is welded to the dash panel member in the close vicinity of the side panel, and a pipe member formed as the instrument panel member has opposite ends welded to the dash panel member in the close vicinity of the side panels and a generally U-shape central portion opening upward and welded to the floor panel at a bottom thereof.\nThe instrument panel member may be provided with brackets welded thereto for mounting a steering shaft and an air-bag on a passengers seat side on an intermediate portion of the instrument panel member between the generally U-shape central portion and one of the opposite ends of the instrument panel member on the passenger seat side. Further, the dash panel member may be provided with a cowl member extending in the transverse direction of the vehicle body and having a closed cross-section.\nThe front body structure of the present invention, including at least one of the dash panel member and the instrument panel member which are formed from a pipe member by hydro-forming, is comprised by an decreased number of parts by spot-welding them at a reduced number of points. This provides the front body with lighter weight and higher structural rigidity as compared with prior art front bodies. In consequent, the front body structure realizes firm and reliable installation of brackets for mounting steering shaft and an air-bag to one of the intermediate portions of the instrument panel member.\nThe front body structure of the present invention, including the dash panel member formed from a pipe member which as the end portion extending downward along the door opening of the side panels and welded to the side sill, provides a wide area for welding. This wide area of welding provides the front body with high structural rigidity.\nThe front body structure of the present invention, including the dash panel member and the instrument panel member formed from pipe members, respectively, which are welded to the inner wall of the vehicle body, improves the structural rigidity of the vehicle body and provides a light weighted front body due to a reduced number of parts.\nFurthermore, the front body structure of the present invention, including the dash panel member formed from a pipe member ensures a sufficient structural rigidity of the front body. This increases the degree of designing a cowl structure because the dash panel member formed from a pipe member provides the front body with a high rigidity."}
{"text": "The present invention relates to a flat knitting machine wherein there is disposed a transfer lock which permits needles to be selectively guided into a loop transfer path or a loop receiving path, in same phase as a knitting lock which permits needles to be selectively guided to three positions, knit, tucking, and welt, so that needles are allowed to knit, tuck, welt or transfer within an operative zone in same phase as the knitting lock in the widthwise direction of a carriage.\nA prior art disclosure relevant to such type of machine is found in Published UK Patent Application GB No. 2080344A, which discloses a cam support comprising at least one arrangement of lifting cams comprising two lifting cams arranged immediately one after the other and each comprising two opposed ramps and a third lifting cam disposed between the adjacent ramps of the two first cams, an empty space being provided below the third cam, said space being connected to two channels limited on one side by the adjacent ramps of the two first cams, so that a pin, respectively of a needle, selected before the arrival of an outer end of said first cams, is entrained by one of the said first cams, in the direction of movement of the cam-support and by the third cam and operates for shaping, whereas a pin, respectively of a needle, selected when said empty space arrives approximately opposite said pin, respectively of a needle, is not entrained by the third cam and operates for picking up.\nIn this prior-art system, however, the cam for pushing a front butt (which is disposed on a needle, while a rear butt is disposed on a pin or jack) to put the needle into a path of transferring operation is movable, and its upper surface constitutes a cam face which will act on the back of the front butt to lift the needle to the transferring path, whereas its underside constitutes a cam face for pushing the front of the front butt to lower the needle into a stitch receiving path. And for the purpose of lowering the needle from the transferring path, there is disposed above said movable cam a guide cam which is fixed. This guide cam acts on the front of the front butt to lower the needle.\nTherefore, to the extent that such fixed guide cam for pushing the abutment surface of the front butt to lower the needle from the transfer path is disposed above said movable cam for putting the needle into the transfer path, the prior-art cam system involves larger lengths of cam plate and needle bed in the longitudinal direction of the needle trick, and also an increase in the length of needle tricks."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of data processing. More specifically, the present invention relates to the employment of namespaces to facilitate local and remote function invocation.\n2. Background Information\nEver since the invention of the first computer, computer scientists have continuously tried to improve the productivity of programmers, such that more applications can be developed with fewer resources to take advantage of the continuous advancements being made in the art of computer and related technologies. First assembler languages were developed to replace machine languages. Then, high level languages, such as FORTRAN, COBOL, PL/I and so forth, were developed to further improve the productivity of programmers. Development of high level languages were followed by structured languages such as Pascal and C, and then object oriented programming languages such as C++. To facilitate development of the Internet and the World Wide Web, “new” languages such as the Hypertext Markup Language (HTML), Java, Javascript, Perl and CGI were developed.\nMost languages also support function calls, which may be user-defined functions or pre-packaged functions, such as those included with the languages' runtime libraries. However, in the earlier days, to be invocable, the functions must exist on the same system as the callers of the functions. Later on, to facilitate development of client-server computing, remote procedure call (RPC), a messaging protocol, was developed to facilitate a program executing on one computer to remotely invoke and access the service of a function on another computer. More recently, to facilitate development of web based applications, and employment of object oriented programming techniques, remote method invocation (RMI) was developed to facilitate interacting with objects on a server.\nWith the advance of public networks, such as the Internet, numerous function resources exist in different parts of the networks, written in different languages on different platforms. It is desirable to be able to invoke these function resources, regardless of their location, host language or host platform without limitation to the relationship between the invoking computing device and the function hosting computing device. Thus, an approach that can efficiently realize this potential is desired."}
{"text": "1. Field of the Invention\nThe invention relates to a roll tape measure, consisting of a housing upper part and housing lower part, a perforation in the housing upper part for a release button, a spring drum with tension spring placed in the housing, whose end is connected with the housing lower part, a back stop for the tension spring, a flexible tape measure rolled up on the spring drum, and a release catch that can be actuated by the release button.\n2. Related Technology\nA roll tape measure of this type can be seen, for example, in DE utility model 69 25 460. In these structural shapes, the release button is connected or integrated with an actuating pin which, on its end in the housing, acts on a release element with a pawl. In this case, the release button must exhibit a smoothly and cleanly finished surface, so that it is guided freely movable in the axis in its lengthwise direction. To achieve this, the release button is designed integral with its end-side stop device as a mushroom-shaped control element. But the release button projects in this case from the housing surface, so that malfunctions may still occur.\nA second roll tape measure of this type can be seen in DE-PS 25 45 203, in which both the housing upper part and the housing lower part exhibits surface parts which are cut out from the housing surfaces and connected by remaining crosspieces with the respective housing part. The surface parts are connected with one another by a pin centrally placed in the housing axis, so that the restoring forces of the crosspieces act in opposite directions. The centrally placed surface part cut out in the housing upper part is designed as a release button, while the surface cut out in the housing lower parts acts as a release catch, which is covered toward the outside by a plate. Relative to the above-indicated known roll tape measure, the release button no longer projects from the housing surface, but forms a relatively smooth surface. The number of individual parts used is relatively small, but it is necessary to provide both housing parts with indentations and to design the connecting crosspieces elastically. To achieve the elasticity of the connecting crosspieces, which are an integral component of the housing, a special selection of the plastics to be used is necessary, so that here only a relatively limited leeway relative to the plastic materials to be used is provided. Further, the cut-out surface parts exhibit relatively long gaps, by which dust and dirt particles can penetrate the interior of the housing.\nIn practice, it has further been shown that a user of such a roll tape measure does not immediately recognize the central arrangement of the release button as such and incorrectly considers the latter to be the axis of the system."}
{"text": "Many computer programs use graphical user interfaces, referred to as GUIs, as an intuitive and convenient means to present information and functionality for users to interact with the computer program. A user can interact with a user interface through a myriad of user interface elements or widgets. A graphical user interface typically has some elements that display information only, such as labels and images, and other elements that receive input from a user through an input mechanism, such as menus, buttons and check boxes. Furthermore, these user interface elements can include a graphical element with a collection of input elements. For example, a tree widget may provide a hierarchy of widgets, where each widget has its own window, dialog or form for interacting with the computer program. In another example, a set of property sheets may have multiple tabs with each tab handling a portion of user input for the user interface. The display and input elements can be placed and arranged within a window, frame, property sheet or other container widget, which can be moved, resized, closed, etc. Users may interact with user interface elements to view display elements and modify input elements across multiple forms or screens.\nThe design and development of graphical user interfaces can be quite complex in considering the myriad of user interface components, and implementing event-driven and parallelism programming paradigms with state and control considerations. Although an event based model may follow a parallelism paradigm, the user interface also needs to have structure and data to track and manage state and control of information in accordance with the desired functionality. For example, it may be desired to have the execution of instructions related to an event behave differently based on the state of user interface components, or the currently assigned values to user interface components.\nFurthermore, the design and development of a graphical user interface can increase in complexity with the increase in the number of user interface elements, the number of generated events, and the complexity of the desired functionality along with the related state and control of the user interface functionality. Moreover, some software programmers may be challenged by the complexities of implementing a design and development of a graphical user interface considering event-driven and parallelism paradigms."}
{"text": "Medical instruments often need to be temporarily connected to peripheral devices and components in the course of operation. An example may be a sampling tube connected to an analyzing instrument such as a capnograph. Another example may an ultrasound probe connected to a sonographic imaging instrument. Such peripheral devices may need to be replaced frequently due to one or more reasons. For example, a disposable probe may be used for each treated patient, and should be replaced after use by a new probe for a next patient. Another reason for frequently connecting and disconnecting probes from an instrument may be related to multi-purpose instruments. Such instruments are configured to carry our one of several routines, for obtaining one of several optional purposes. Generally, a particular routine and purpose may be associated with a specific peripheral device that needs to be connected to the instrument for carrying out the routine. For example a user of a sonographic imaging system may wish to carry out one type of examination using one type of ultrasound probe, and then to carry out a second type of examination using a second type of ultrasound probe. Thus, frequent replacement of probes is required, typically being carried out by disconnecting a previously used probe and connecting a new probe to the instrument, instead."}
{"text": "Due to the ever increasing demand for data storage in today's digital world, distributed storage systems have received a lot of attention. In distributed storage systems, data as a whole or different portions of data are stored over several physically independent computer storage servers (also known as storage nodes). The total storage capacity of the distributed storage system is designed to be much greater than that of a single storage node.\nThere is an exponentially increasing amount of data produced and stored worldwide in distributed storage systems. A large fraction of that data is duplicated, for instance, as a result of unintentional copies of all or part of the data being stored multiple times, therefore leading to wasted storage resources. In particular, company/enterprise data, consisting mostly of structured formats such as documents and spreadsheets, is often replete with unnecessary duplicate data sometimes estimated over ˜70%. Therefore, an enterprise-wide storage system can make significant space savings if the data can be de-duplicated and stored only once. However, in many distributed database systems, deduplication is difficult to implement as the data models supported by, for example, relational database management systems (RDBMS), are not designed with deduplication in mind.\nFile-system level deduplication products are alternatives that aim to work with raw data on disks underneath any higher level storage abstractions such as databases, thereby leaving the upper or higher layers to operate as normal (e.g., unaware of the deduplication process underneath it). However, it has been shown that file-system level deduplication does not lead to optimal storage resource savings (e.g., for databases, especially) as there are too many small changes (e.g., file headers, timestamps, etc.) that make blocks of data subtly different and thereby defeat deduplication efforts. File-system level deduplication can also introduce additional data loss risks, as the database on top can make incorrect assumptions about which critical data could be lost due to a physical disk error.\nIn view of the above, there is a need for efficient systems and methods for data storage that can overcome these challenges and store the data in a distributed storage system securely, reliably and with minimal duplication."}
{"text": "In the above-noted co-pending application, there is disclosed a method of co-nucleating precious metal powders with a small percentage of a refractory oxide, most particularly titania or zirconia. The process involves a slow precipitation from an acidic solution with ammoniacal hydrazine, followed by a heat treatment. Such powders exhibit increased sintering temperatures, and find utility when formulated as pastes, with a small amount of glass frit as binder, in thick film electronic applications, and in catalysis.\nThe present invention is based, in one part, on the discovery that with a somewhat altered procedure, the process described above can be applied to nickel and the product exhibits a significantly increased oxidation resistance. The present invention is based, in a second part, on the discovery that with an essentially different process, precious metal powders of both increased density and higher sintering temperatures are produced.\nIn the above-mentioned co-pending application, it is noted that the addition of refractory oxides as distinct ingredients in certain thick film compositions is known, that the addition of copper oxide to gold to make a so-called \"fritless\" thick film conductors is known, and that it is also known to improve bulk metal properties such as creep strength and tensile strength with a refractory oxide as a distinct second phase, e.g. dispersion hardening. The surprising and novel aspect of the present invention, as well as that described in said co-pending application, is the significantly improved thermal properties of metal powders produced in accordance herewith."}
{"text": "The separation of light gases has become increasingly important over the years particularly in areas relating to natural gas recovery and enhanced oil recovery where there can be large amounts of light gases present. Often times, there are appreciable amounts of nitrogen present, as well as other smaller molecules such as carbon dioxide, oxygen and water, which need to be separated from other larger molecules such as methane, ethane, propane, butane and corresponding olefins.\nIn the natural gas industry, the separation of nitrogen from hydrocarbons, and particularly methane, is of primary interest in order to avoid nitrogen contamination from gas wells into natural gas pipelines. Natural gases which contain significant amounts of nitrogen may not meet minimum heating value specifications, reduce pipeline capacities and require additional compression horsepower and fuel consumption. Furthermore, as high quality gas reserves are gradually being depleted, the natural gas industry is becoming more dependent on lower qualtiy gas which contains larger amounts of nitrogen and other undesired gases. Accordingly, nitrogen removal from natural gases has attained increased importance.\nThe depletion of oil reserves has also influenced the petroleum industry where successful recovery of petroleum often requires the use of an enhanced recovery technique. One such often used technique involves the injection into the reservoir of a fluid which will not support combustion; an often used fluid for this technique is nitrogen or a nitrogen-containing gas. However, the use of this technique increases the level of nitrogen contaminant in the gas fraction recovered from the reservoir, i.e., the petroleum gases, above their naturally-occurring nitrogen concentration.\nNitrogen injection for enhanced oil recovery introduces a further problem because the nitrogen concentration in the petroleum gases, also known as case-head or well-head gases, does not remain constant over the life of the recovery operation. Although the nitrogen concentration variation will strongly depend upon particular reservoir characteristics, a general pattern is predictable. Typically during the first few years that enhanced recovery with nitrogen injection is employed, the nitrogen concentration in the gases may remain at about the naturally-occurring level, increasing thereafter as more nitrogen injection is required to recover the oil.\nHence, in order to render the use of nitrogen effective for enhanced oil and gas recovery, processes which are both cost effective and tolerant to compositional variations are required to separate nitrogen from methane and other gases.\nThe nitrogen contaminant can be removed from the nitrogen-containing gas by distillation. For example, U.S. Pat. No. 4,415,345 discloses a solution to this problem by providing processes that utilize distillation to remove nitrogen from natural gas streams, wherein a nitrogen heat pump is employed with both single and double distillation column arrangements to process streams of varying nitrogen content. In general, however, distillation processes such as described in the above-mentioned patent are well suited only for larger-scale operations due to their relatively high equipment costs and complex operating procedures.\nIn smaller-scale natural gas operations as well as in other areas such as synthesis gas and coke oven gas processing, adsorption processes can be especially well suited. The adsorption capacities of adsorption units can, in many cases, be readily adapted to process gas mixtures of varying nitrogen content without equipment modifications, i.e., by adjusting adsorption cycle times. Moreover, adsorption units can be conveniently skid-mounted, thus providing easy mobility between gas processing locations. Further, adsorption processes are often desirable because more than one component can be removed from the gas. As noted above, nitrogen-containing gases often contain other gases that contain molecules having smaller molecular dimensions than nitrogen, e.g., for natural gas, carbon dioxide, oxygen and water, and for coke oven gas, carbon monoxide.\nHowever, despite the advantageous aspects of adsorption processes, the adsorptive separation of nitrogen from methane has been found to be particularly difficult. The primary problem is in finding an adsorbent that has sufficient selectivity for nitrogen over methane in order to provide a commercialy viable process. In general, selectivity is related to polarizability, and methane is more polarizable than nitrogen, i.e. The polarizability volumes for nitrogen and methane are respectively, N.sub.2 =17.6.times.10.sup.-25 cm.sup.-3 ; CH.sub.4 =26.0.times.10.sup.-25 cm.sup.-3. See P. W. Atkins, PHYSICAL CHEMISTRY, p. 773, Freedman (1982). Therefore, the equilibrium adsorption selectivity of essentially all known adsorbents such as large pore zeolites, carbon, silica gel, alumina etc. all favor methane over nitrogen. However, since nitrogen is a smaller molecule than methane, it is possible to have a small pore zeolite which adsorbs nitrogen faster than methane. Clinoptilolite is one of the zeolites which has been disclosed in literature as rate selective adsorbent for the separation of nitrogen and methane.\nThe separation of gaseous mixtures of methane and nitrogen using both raw clinoptilolite and clinoptilolite which had been ion-exchanged with calcium cations is described in the following publication; T. C. Frankiewicz and R. G. Donnelly, METHANE/NITROGEN SEPARATION OVER THE ZEOLITE CLINOPTILOLITE BY SELECTIVE ADSORPTION OF NITROGEN, Chapter 11, INDUSTRIAL GAS SEPARATION, American Chemical Society, 1983. They disclose that at long adsorption times, adsorption approaches thermodynamic equilibrium and there is a tendency for adsorbed nitrogen to be replaced by methane. However, since methane diffusion is slower than nitrogen diffusion into clinoptilolite, the separation can be made on a rate basis.\nJapanese Patent Application (Kokai) No. 61-255,994 discloses a process for producing a high-caloric gas comprising two adsorption zones wherein nitrogen and other non-combustible low-caloric components are removed from a feed gas, e.g., coke oven gas or methane reaction gas, which also contains hydrogen, methane and other hydrocarbons. This Japanese patent application discloses that the nitrogen is adsorbed on a clinoptilolite adsorbent that may be naturally produced clinoptilolite, natural clinoptilolite that has been crushed as required either in its original form or after ion exchange or other chemical treatment, natural clinoptilolite that has been combined with a suitable binder, then compacted and sintered, natural clinoptilolite that has merely been heat-treated, or from clinoptilolite obtained by a synthetic process. This Japanese patent application does not, however, disclose any specific cations that would be suitable as ion-exchange agents in clinoptilolite for adsorbing nitrogen or smaller molecules.\nJapanese Patent Application (Kokai) No. 62-132,542 discloses an adsorbing and separating composition composed of a clinoptilolite type zeolite containing calcium cations in a mole ratio of CaO/Al.sub.2 O.sub.3 of 0.4 to 0.75. The application discloses that the composition is useful for separating molecules with a kinetic diameter of less than 3.7 angstroms from molecules with that of 3.7 angstroms or greater, e.g., removal of small quantity of nitrogen from methane gas, or bulk separation of nitrogen from a methane-containing coke oven gas or coal mine draught gas, etc.\nAlthough the above-described references propose various adsorbents and processes for separating nitrogen from methane other gases, more efficient and effective adsorbents and commerically viable processes are sought. In particular, it is desired to provide adsorbents and processes that are particularly useful for separating nitrogen or smaller molecules from methane or larger molecules."}
{"text": "The present invention generally relates to an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture. More specifically, the present invention relates to an apparatus, a system and a method for securing, for attaching and/or for detaching a portable electronic device, such as, for example, a camera, a camcorder, a digital camera, a digital music player, a video game, a digital video player, a digital video recorder, a personal data assistant (hereinafter “a PDA”), a laptop or tablet computer, an E-book, an E-reader, a cellular telephone and/or the like. The apparatus, the system and the method for securing, for attaching and/or for detaching the device may attach the device to the fixture to deter and/or prevent theft and/or destruction of the device. The device may be a portable electronic device which may be, for example, displayed for sale by a retailer, a wholesaler and/or the like.\nIt is generally known, for example, that vendors, retailers and/or wholesalers may display a device to a customer at, for example, a retail store and/or a sales facility for examination and/or inspection by the customer. The device may be a portable electronic device, such as, for example, a camera, a digital camera, a compact disc player, an MP3 player, a PDA, a laptop or tablet computer, an E-book, an E-reader, a cellular telephone and/or the like. The device is traditionally displayed in conjunction with a fixture, such as, for example, a wall, a floor, a pillar, a support beam, a stair case or a fixture, such as, for example, a cabinet, a table, a shelf and/or the like.\nGenerally, the device is fixed to a surface of the fixture, and/or a cable attaches the device to the fixture. The cable is often not adjustable or has a limited length. In addition, the cable is often not detachable. As a result, the customer and/or the user is unable to examine the device while being independent of the cable and/or fixture. For example, pulling of the cable often prevents the customer and/or the user from examining the various characteristics of the device, such as, for example, the weight, the feel, and the portability of the device. Further, the attachment of the cable prevents the customer and/or the user from examining the size, weight and/or other dimensions of the device with respect to a laptop bag, a book bag, and/or a pant pocket, for example.\nA need, therefore, exists for an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture. Additionally, a need exists for an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture which may allow a device to be detached from the fixture to allow a customer and/or a user to manipulate and/or examine the device. Moreover, a need exists for an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture which may allow a device to be re-attached to a fixture after detachment from the fixture. Further, a need exists for an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture which may provide electrical security to the device while the device is detached from a fixture. Even further, a need exists for an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture which may utilize a retractable reel for retracting a cable. Moreover, a need exists for an apparatus, a system and a method for securing, for attaching and/or for detaching a device to a fixture which may trigger an alarm for a separation of the device from the cable and/or the separation of the device from the apparatus."}
{"text": "The invention relates to a high-voltage apparatus. Furthermore, the invention relates to a corresponding external reproduction apparatus. Furthermore, the invention relates to a corresponding system having the high-voltage apparatus and the external reproduction apparatus.\nIf an accident occurs in or with a high-voltage apparatus, it is often impossible for people outside the high-voltage apparatus to tell whether the high-voltage apparatus is live and whether a threat is therefore posed by the high-voltage apparatus.\nThe object on which the invention is based is that of providing a high-voltage apparatus and an external reproduction apparatus that helps to make it possible to ascertain without touching whether a threat is posed by the high-voltage apparatus.\nThis and other objects are achieved by a high-voltage apparatus having at least one high-voltage component, at least one current sensor that is associated with a respective high-voltage component, and at least one transmission apparatus that has at least one respective associated current sensor, wherein the at least one transmission apparatus is designed to take a measurement signal from the at least one respective current sensor as a basis for providing an information signal via a wireless radio link, which information signal has a piece of information about whether or not a potential threat is posed by the respective high-voltage component, for the purpose of reproducing the piece of information on an external reproduction apparatus.\nThe invention is distinguished by a high-voltage apparatus. The invention is furthermore distinguished by a corresponding external reproduction apparatus. The invention is furthermore distinguished by a system having the high-voltage apparatus and the external reproduction apparatus.\nThe high-voltage apparatus has at least one high-voltage component. The high-voltage apparatus has at least one current sensor that is associated with a respective high-voltage component. The high-voltage apparatus has at least one transmission apparatus that has at least one respective associated current sensor. The at least one transmission apparatus is designed to take a measurement signal from the at least one respective current sensor as a basis for providing an information signal by way of a wireless radio link. The information signal has a piece of information about whether or not a potential threat is posed by the respective high-voltage component. The at least one transmission apparatus is designed to provide the information signal by way of a wireless radio link.\nBy way of example, for each current sensor, the high-voltage apparatus has a transmission apparatus that has only one respective associated current sensor.\nExamples of such high-voltage components are components that are provided for producing, transmitting and/or consuming electric power at a high electrical voltage. A high electrical voltage of this kind is at least 300 V for example. A high-voltage component of this kind is an electrical energy store, an electric drive machine, a cable, a plug and/or a connection in the high-voltage apparatus, for example.\nThe current sensor is installed in a manner integrated in a plug and/or in what is known as a black box and/or in a cable, for example.\nAn example of a high-voltage apparatus of this kind is a vehicle, such as a hybrid or electric vehicle, a test bench, a photovoltaic installation, a power supply installation, a rail vehicle installation and/or a rail vehicle.\nA high-voltage apparatus of this kind allows the information signal to be provided for the external reproduction apparatus, which information signal has a piece of information about whether or not a potential threat is posed by the respective high-voltage component. By way of example, the information includes whether or not the respective at least one current sensor measures an electric current and/or what the level of the measured electric current is and/or whether the respective at least one current sensor has a fault. This means that it is possible for an emergency worker, for example, to establish whether a threat is posed by the high-voltage apparatus and whether or not he or she can approach the high-voltage apparatus without danger in a simple manner and without touching.\nAccording to a further advantageous embodiment, the information signal has a piece of position information that is representative of a position of the at least one respective current sensor in the high-voltage apparatus.\nThis means that it is a very simple matter for the emergency worker, for example, to establish where the respective current sensor is positioned in the high-voltage apparatus. It is therefore possible to establish what location in the high-voltage apparatus may pose a threat.\nAccording to a further advantageous embodiment, the high-voltage apparatus is a vehicle.\nSpecifically in the case of a vehicle, such as a hybrid vehicle and/or an electric vehicle, it is important for the emergency worker to know whether or not he or she can approach the vehicle in order to be able to initiate rescue measures quickly and safely.\nAccording to a further advantageous embodiment, the information signal has a vehicle identifier that is representative of the vehicle type of the vehicle.\nSpecifically in the case of modern vehicles, it is often not immediately evident from the outside whether an electric vehicle and/or a hybrid vehicle is involved. Transmission of the vehicle type allows the rescue worker, for example, to tell whether the vehicle is an electric vehicle and/or a hybrid vehicle.\nThe external reproduction apparatus has a radio reception unit that is designed to receive the information signal that has been provided by the high-voltage apparatus or by an advantageous embodiment of the high-voltage apparatus. The external reproduction apparatus is designed to react to reception of the information signal by displaying a respective two-colored set of traffic lights that is associated with the at least one respective current sensor. A first traffic light color is representative of no threat being posed by the respective high-voltage component. A second traffic light color is representative of a threat being posed by the respective high-voltage component.\nThe external reproduction apparatus has a display, for example, such as a screen, for reproduction. The external reproduction apparatus is particularly a portable apparatus, such as a smart phone and/or a tablet PC.\nThe respective two-colored set of traffic lights provides a simple way of using the external reproduction apparatus to reproduce whether or not the respective high-voltage component poses a threat. By way of example, the first traffic light color may thus be green and the second traffic light color may be red.\nAccording to one advantageous embodiment, the information signal has a piece of position information that is representative of a position of the at least one respective current sensor in the high-voltage apparatus. The external reproduction apparatus is designed to display the respective two-colored set of traffic lights at the position of the at least one respective current sensor.\nThis means that the position of the at least one respective current sensor can be displayed on the external reproduction apparatus in a simple manner. By way of example, it is thus a simple matter for the rescue worker to establish where the respective current sensor is positioned in the high-voltage apparatus. Hence, it is possible to establish what location in the high-voltage apparatus may pose a threat.\nAccording to a further advantageous embodiment, the high-voltage apparatus is a vehicle. The information signal has a vehicle identifier that is representative of a vehicle type of the vehicle. The external reproduction apparatus is designed to display the respective two-colored set of traffic lights at the position of the at least one respective current sensor on a vehicle depiction that matches the vehicle type of the vehicle.\nThis means that the external reproduction apparatus can be used to display simply and clearly where the respective at least one current sensor and the respective high-voltage component are arranged in the vehicle.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of one or more preferred embodiments when considered in conjunction with the accompanying drawings."}
{"text": "1. Field of the Invention\nThis invention relates to the field of electric guitars of the Spanish type, as distinguished from Hawaiian or steel guitars.\n2. Description of Prior Art\nIt has long been conventional, in electric Spanish guitars, to stretch the strings over a bridge on the body of the guitar and over a nut located at the junction between the fingerboard and the head. Thus, the long operative or \"working\" portions of the strings are defined between the bridge and the nut. After passing over the nut, in notches or narrow grooves therein, the strings connect to tuning screws located on the head.\nAt least relative to two of the strings on one major type of electric Spanish guitar, there are \"string guides\" on the head, beneath which such strings pass after bending over the nut. The string guides are actually hold-down elements to make sure that the strings do not become even momentarily disengaged from the nut when the amplitudes of vibration are great.\nThe increasing of the dwell time of a vibrating guitar string is a major goal of guitar designers, since musicians want the sound produced by each plucking action to continue for a substantial time period. Stated conversely, it is desired that the decay rate (decrease in amplitude of vibration) be low. Prior-art attempts to increase dwell (or \"sustain\") and reduce decay have been unsuccessful to date, particularly because in Spanish guitars (as distinguished from conventional Hawaiian guitars), the thin, long wooden necks tend to vibrate (resonate) in response to string vibration and this represents lost mechanical energy which reduces dwell. For the same pitch, in the open-string condition, the string of a Spanish guitar is conventionally smaller in diameter--and under less tension--than is the string of a Hawaiian guitar.\nIt is possible to increase dwell substantially, even in Spanish guitars, by using heavy string-anchoring elements, unsightly elements, complex and expensive elements, etc. However, in no case known to applicant has the dwell been increased markedly by a simple, inexpensive, light, good-looking element which need have no movable parts. The present invention does achieve this long-desired result and, furthermore, achieves the advantage of making it possible to eliminate the above-indicated string guides with attendant cost saving."}
{"text": "1. Field of the Invention\nThe invention relates to a projection illumination system.\nLithographic methods are conventionally employed for the manufacturing of miniaturized structures and components. Miniaturized components and structures include, for example, integrated circuits, fluid crystal elements, micromechanical components and the like. Hereby, patterns or structures predetermined on a mask (“reticle”) are imaged onto a radiation sensitive substrate, such as, for example, a wafer provided with a radiation sensitive layer (“resist”), by an imaging optics of a projection illumination system. The light required for the imaging is generated by an illumination system of the projection illumination system, which comprises an illumination optics for illuminating the mask.\nWith the continued trend towards increasing miniaturizing of the components to be manufactured, greater and greater demands are made on the imaging quality of projection illumination systems.\n2. Brief Description of Related Art\nA projection illumination system is known, for example, from US 2002/0001088 A1, in which particular parameters characterizing the imaging quality can be determined by an interferometric measurement of the imaging optics. An imaging optics is, for example, known from U.S. Pat. No. 5,973,863 which comprises actuators which adjust a position of at least one optical element of the imaging system relative to other elements of the imaging system. Thus it is possible to reduce certain imaging errors of the imaging optics by correspondingly controlling the actuators.\nDue to the large number of optical components in a projection illumination system, the imaging quality is determined by a multitude of parameters of the optical components and is also subject to operation related drifts."}
{"text": "The Internet Engineering Task Force (IETF) has defined protocols in the area of Layer Two Tunneling Protocol—Version 3 (L2TPv3) and Pseudo Wire Emulation Edge-to-Edge (PWE3) Architecture, as defined in RFC 3193, RFC 3931, RFC 3985, RFC 4591, and RFC 4719. The L2TPv3 pseudo wires may support multiple types of layer two encapsulations or “Pseudowire Types”, including protocols similar in packet format to IEEE 802.3, virtual LAN tagging protocols such as IEEE 802.1Q, IEEE 802.1ad, and IEEE 802.1ah, as well as Point-to-Point Protocol (PPP), Frame Relay, and other layer 2 technologies.\nAs defined in the standard track protocols above, the Pseudowire Types including Ethernet, VLAN, PPP, or Frame Relay, etc are encapsulated into “L2TPv3” to emulate a single layer 2 tunnel link between exactly two endpoints.\nEnabling L2TPv3 on cable modems (CM), embedded multimedia terminal adapters (E-MTA), cable modem routers (CM-R), cable modem with Circuit Emulation Service (CES), cable modem with IPSec or Firewall and/or similar DOCSIS Customer Premise Equipment (D-CPE) will allow L2TPv3 IP packets to be encapsulated into Data Over Cable Service Interface Specification (DOCSIS) data frames for two-way data transmission in a cable system. Devices with this capability are referred to herein as “L2TPv3 enabled DOCSIS CPE”, or simply “Enhanced CPE”."}
{"text": "This invention relates to a cam shaft drive arrangement for a four-cycle engine and more particularly to such a drive arrangement for a four-cycle engine of an outboard motor.\nIn many four-cycle engines, the cam shaft is driven from the crankshaft by means of a flexible transmitter in the form of a toothed belt. When belt drive cam shaft timing arrangements are employed, the belt is positioned outside of the actual engine body. That is, unlike gear or chain timing mechanisms, the belt should be protected from the engine lubricant and preferably positioned externally of the main engine body. This presents some particular potential difficulties.\nThat is, because the timing drive belt is positioned externally of the engine, this means that the drive for the cam shaft is spaced from the outermost bearing that supports both the crankshaft and the cam shaft. Thus, bending loads are placed on the cam shaft that can cause high wear to either the bearing and/or the shaft which is borne.\nIt is, therefore, a principal object of this invention to provide an improved timing belt arrangement for four-cycle internal combustion engines.\nIt is a still further object of this invention to provide an improved bearing and belt drive arrangement for the cam drive of a four-cycle engine.\nThe aforenoted problems are particularly acute in conjunction with outboard motors. As is well known, such outboard motors are quite compact in nature. Because of the premium for space, the positioning of the cam shaft drive on the end of the engine, as is necessitated by the use of a drive belt, tends to increase the overall length of the engine. This is particularly undesirable with applications where space is at a premium, such as with an outboard motor application.\nThese problems are further aggravated when other accessories are driven off of the engine. That is, these further accessory drives can add considerably to the length of the engine.\nIt is, therefore, a still further object of this invention to provide an improved cam shaft drive arrangement employing a drive belt for use in outboard motors."}
{"text": "1. Field of the Invention\nThis invention relates to webbing clamps and, more particularly, to a novel webbing clamp and resilient hoop apparatus and method for releasably mounting an airspeed indicator to webbing of a paraglider harness, the resilient hoop being secured to the webbing clamp by a unique ball joint.\n2. The Prior Art\nWithin the past decade or so there has been a significant growth in the sport of paragliding. The name itself is a combination of two words parachute (para) and hang gliding (gliding) and, as the name implies, constitutes the unique sport of manned flight in reliance on the principles of flight embodied in these two sports. Instead of a conventional parachute canopy having a generally hemispherical profile, the paraglider is configured with a wing-like structure from which the pilot is suspended by the shroud lines and webbing support system of the harness. The paraglider wing is configured as an airfoil having individual cells extending transversely across the wing. Each cell includes an opening in the front to allow ram air resulting from the forward motion of the flight to inflate each cell and thereby create the airfoil profile of the paraglider wing.\nThe pilot in his/her harness is suspended below the paraglider wing by the shroud lines so that there is no rigid structural framework in the paraglider. This means that any airspeed indicator (such as disclosed in my copending patent application Ser. No. 08/065,989 filed May 24, 1993 for AIRSPEED INDICATOR APPARATUS AND METHOD) will be difficult to affix to the paraglider. This is particularly relevant since the webbing of the paraglider harness does not extend from the pilot along normal horizontal and vertical axes. Specifically, the webbing radiates outwardly and upwardly from the pilot at various angles so that even the plane of the webbing itself will be at an angular offset from the foregoing axes. Further, since it is essential that my airspeed indicator is oriented in the proper direction in order to provide an accurate indication of airspeed, it is equally essential that the airspeed indicator is adjustably mountable to the webbing of the paraglider harness. Not only must the airspeed indicator be adjustably secured to the webbing of the paraglider harness, it must also be releasable in the event a shroud line inadvertently becomes snagged across the airspeed indicator. This requirement places a fairly rigorous set of preconditions on any system used to secure an airspeed indicator to the webbing of a paraglider harness.\nIn view of the foregoing, it would be an advancement in the art to provide a webbing clamp for use in removably, adjustably, and releasably securing an airspeed indicator to the webbing of a paraglider harness. It would also be an advancement in the art to provide an engagement mechanism for releasably engaging the barrel of an airspeed indicator, the engagement mechanism being formed as a hoop of resilient plastic to accommodate the hoop resiliently encircling the barrel of the airspeed indicator. Another advancement in the art would be to provide a ball and socket engagement mechanism between the resilient hoop and the webbing clamp. Such a novel apparatus and method is disclosed and claimed herein."}
{"text": "This invention relates, generally, to a fireplace tool or accessory and, more specifically, to a fireplace tool for flue damper system viewing, opening and closing.\nA flue damper system is generally provided in a working fireplace for opening and closing the flue of the fireplace and, perhaps, controlling the draft of air through the flue for a fire. The damper system is generally located in the ceiling of the fire chamber of the fireplace with a downwardly-depending lever for manually adjusting the damper.\nAccordingly, the lever is covered with soot from previous fires therebelow and can burn a person's hand if touched when fire heated. Many persons therefore try to use a poker for moving the damper, but its straight shape and the location of the damper system makes this awkward. Also, the location of the damper and even the depending damper-adjusting lever generally does not permit a person to see either directly without getting his (her) head near the floor or inside the fire chamber. Adjustments are often incorrect, therefore, leading to a fire with a closed damper and resulting smoke-filled room or no fire with an open damper and resulting draught. This situation is, therefore, in need of improvement."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for coating an object with a photosensitive material, and more particularly to a method and an apparatus for uniformly coating an object with a photosensitive material.\n2. Description of the Related Art\nAs a conventional method for coating a semiconductor substrate with a photosensitive material in a sheet form, a spin coating method is used. In the spin coating method, the photosensitive material is dropped and coated on a substrate that is being quickly rotated. That is, the photosensitive material is coated by a centrifugal force that is applied to the substrate. The spin coating method is proper to coat a small substrate such as a wafer with the photosensitive material. However, the spin coating method may be improper to coat a large and heavy substrate such as an LCD panel with the photosensitive material.\nIn order to coat a large and heavy substrate with a photosensitive material, a slit coating method has been proposed. In the slit coating method, the photosensitive material is supplied through a slit nozzle having a length thereof longer than a width thereof and is coated in a sheet form.\nFIG. 1 is a schematic view showing a conventional slit coater.\nReferring to FIG. 1, a slit coater 10 changes a photosensitive material having an indeterminate form into a photosensitive material having a band form, which has a length thereof longer than a width thereof. The slit coater 10 is aligned at a designated position “A” so as to supply the photosensitive material onto a substrate 40.\nFIGS. 2A and 2B are views showing a silt coater and a substrate partially coated with photosensitive material.\nReferring to FIG. 2A, the slit coater 10 supplies the photosensitive material 20 from a photosensitive material supply apparatus 30 to the substrate 40. The photosensitive material 20 supplied from the slit coater 10 has a band form as shown in FIG. 2B. The slit coater 10 coats the substrate 40 with the photosensitive material 20 in a first direction and at a first speed.\nFIGS. 3A and 3B are views showing the substrate totally coated photosensitive material by the slit coater.\nReferring to FIGS. 3A and 3B, the slit coater 10 coats the substrate 40 with the photosensitive material 20 from the first position “A” to a second position “B”. In the slit coating method, the photosensitive material may be easily coated on an area of substrate greater than that of a wafer.\nFIG. 4 is a sectional view showing a profile of the photosensitive material coated on the substrate according to the slit coating method.\nReferring FIG. 4, in the slit coating method, a thickness T1 at the first position “A” and a thickness T2 at the second position “B” are different from a thickness T3 at a center of the substrate 40, which is unpreferable. Also, the slit coating method causes coating failure as wave patterns or scratch patterns in the coated photosensitive material. Accordingly, the slit coating and spin coating methods are successively performed so as to coat the photosensitive material on the substrate.\nWhen both the slit coating method and the spin coating method are performed on the same substrate, the time for coating the substrate with the photosensitive material may greatly increase. Also, since equipments for performing the slit coating method and the spin coating method are respectively required, cost of equipments may increase."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical transceiver, in particular, the invention relates to an optical transceiver with an enhanced electro-magnetic interference (hereafter denoted as EMI) tolerance.\n2. Background Art\nAn optical transceiver generally includes a transmitter optical subassembly (hereafter denoted as TOSA) to transmit signal light, a receiver optical subassembly (hereafter denoted as ROSA) to receive another signal light, an electronic circuit communicating with the TOSA and the ROSA, and a housing to enclose the TOSA, the ROSA and the circuit therein. One type of the optical transceiver is called as the pluggable optical transceiver in which the optical transceiver may be plugged within a cage prepared in the host system without turning the electrical power of the system off.\nFIG. 10 illustrates such a pluggable optical transceiver 100. The optical transceiver 100 shown in FIG. 10 is going to be set within the cage 102 on the circuit board 104 of the host system through the port 103a prepared in the face panel 103 of the system. The optical transceiver 100 provides an optical receptacle 106 in the front end thereof, into which an external optical connector 107 is set to secure the optical coupling between the external optical fiber in the optical connector 107 and active devices installed in the TOSA and the ROSA. The cage 102 provides an electrical connector in the deep end thereof. Setting the optical transceiver 100 within the cage, the electrical plug formed in the rear end of the optical transceiver 100 mates with the electrical connector in the cage to establish the communication path between the optical transceiver 100 and the host system.\nThe EMI radiation from electronic equipment has been a subject for a long time. Because the pluggable optical transceiver is used in an arrangement to be inserted within the cage, electro-magnetic interference (hereafter denoted as EMI) radiation may leak through a physical gap between the optical transceiver and the cage. In order to reduce such EMI radiation, the optical transceiver generally provides ground fingers or ground fins that makes in contact with the inner surface of the cage when the optical transceiver 100 is set in the cage.\nA United States patent, the U.S. Pat. No. 7,491,090, has disclosed one type of elastic fingers attached to surround the optical receptacle. When the optical connector is set in the optical receptacle, the optical connector pushes the elastic finer outwardly to make in contact with the cage; while, the optical connector is released from the optical receptacle, the elastic finger is apart from the cage to facilitate the extraction of the optical transceiver from the cage. A Japanese Patent Application published as JP-2003-270492A has disclosed another type of an optical transceiver having a metal shield to fill a gap between the optical receptacle and the port.\nConventional mechanisms to reduce the EMI radiation are those described above; that is, the ground fingers surrounding the optical receptacle not only make in contact with the cage to fill the gap against the cage, or at least to narrower the gap, but also stables the chassis ground or case ground.\nHowever, recent optical transceiver operates in a higher frequency, sometimes exceeding 10 GHz; a small gap causes the leak of the EMI radiation from the optical transceiver. For instance, misalignment between the face panel and the cage may leak the EMI radiation. Moreover, the conventional ground fingers inherently have a gap with respect to adjacent fingers, which may also cause the leak of the EMI radiation with higher frequencies. Further, the external optical connector set in the optical receptacle also forms a small gap against the inner wall of the optical receptacle. Although some optical connectors implement a member to shield such a gap, a small sized connector applied in recent systems often have no room to place such a shielding member.\nThus, a way to implement an auxiliary member to shield a small gap shows a limitation as the operating frequency becomes higher. The present invention is to provide a technique to reduce the EMI radiation from an optical transceiver independent of ways to cover the small gap."}
{"text": "In online gaming, game hosting services and game developers have created a number of ways to track and personalize the online gaming experience. A drawback of existing systems is that many of the features have grown up independent of each other. Games send blobs of data about gamers back and forth to a central service, but the service has no way to understand and aggregate the data outside of the game context. Games can host their own Websites, but the data displayed there is not universally accessible to other games.\nIn a sense, then, the service and games offer two parallel communities that offer great—but separated—resources for gamers. First, in the game community, while playing a game, the gamer can see the community of others who play the specific game, the leaderboards for that game, and his personal achievements in that game. A game can tell a gamer, from the Service data, if a Friend is online, but it can't tell the gamer what, exactly that Friend is doing on the Service or when he will be available.\nSecond, in the service community, the service knows a gamer player's history, all of the games he's played, the amount of time he spends online, the size of his Friends list and all of the games that Friends have played or are playing, and feedback the gamer has given and received.\nSystems have tried to leverage these on-line communities to match various players to allow them to play multi-player games. Nevertheless, in general such systems, which typically emphasize skill or experience in a single game or small family of games, there is a desire to provide mechanisms to find similarly skilled players or to determine the skill level"}
{"text": "1. Field of the Invention\nThe present invention relates to an arrangement of a seat belt anchor in an automotive seat, and in particular to an arrangement for covering the seat belt anchor in the seat.\n2. Description of Prior Art\nReference is made to FIGS. 1 and 2, which shows a typical conventional seat belt anchoring arrangement found in the automotive seat. According thereto, as shown in FIG. 2, an anchor member (11') is mounted on a reclining device bracket (5'). The anchor member (11') is connected to a seat belt (1'). Although not shown clearly, a reclining device is mounted on the bracket (5') such that the upper arm of the reclining device is connected to the lower side of a seat back (SB'), while the lower arm of the same is fixed on the bracket (5'), thereby allowing the seat back (SB') to be adjustably inclined forwardly and backwardly relative to the seat cushion (SC'). The bracket (5') is fixed fast on a slide rail (SR').\nThose bracket (5'), reclining device, and seat belt anchor member (11') are normally covered with a cover member (C) in order to maintain the appearance of the seat aesthetically. The cover member (C) is formed with a cut-away portion (C2) through which the seat belt (1') passes. Also, the cover member (C) is attached laterally of the seat cushion (SC') by means of several securing brackets (C3). Some of the securing brackets (C3) are, as shown, fixed to the extension of an upper rail (SR1') associated with the slide rail (SR').\nHowever, with this conventional arrangement, it has been found that the cut-away portion (C2) is defined relatively wider, and due to the interposition of such securing brackets (C3), there is a spacing between the inward edges of the cover member (C) and the lateral wall (SC1') of seat cushion (SC'). This hitherto construction leads to the drawbacks that (i) dust or small foreign materials or particles are dropped through those cut-away portion (C2) and spacing, and accumulated within the cover member (C) as well as upon the reclining device and seat belt anchor member (11'), which lowers the durability and operation of the reclining device and seat belt anchor (11'), and (ii) the cut-away portion (C2) of cover member (C) is still objectionable, in that the cut-away portion (C2) per se presents an uneven irregularity or poor design effect from the aesthetic viewpoint and is thus insufficient to improve the outward aesthetic appearance of the cover member (C), especially in the spot where the seat belt (1') crosses or extends outwardly from the cut-away portion (C2), and additionally the interior within the cover member (C) can easily be viewed through that cut-away portion (22)."}
{"text": "Technical Field\nThe present invention relates to a rotary joint.\nBackground Art\nA rotary joint is used in relative rotation portions of various rotary apparatuses such as a semiconductor manufacturing apparatus, a medical apparatus, and an apparatus for food. The rotary joint includes a first member and a second member that rotate relative to each other, and it is necessary to cause a plurality of different types of fluids between these two members to flow through different passages. Thus, the rotary joint has a configuration in which a plurality of first flow passages are formed in the first member, a plurality of second flow passages are formed in the second member, and one first flow passage and one second flow passage are connected to each other in a one-to-one relation, for causing the fluids to flow therethrough (see, e.g., Japanese Laid-Open Patent Publication No. 2002-22076).\nFIG. 4 is a cross-sectional view of a conventional rotary joint. The rotary joint 90 includes a cylindrical case body 91 (the above first member) and a shaft body 92 (the above second member) provided inside the case body 91. Upper and lower two first flow passages 95 are formed in the case body 91, and upper and lower two second flow passages 96 are formed in the shaft body 92. An opening of the upper first flow passage 95 and an opening of the upper second flow passage 96 are located at the same position in an axial direction, and these flow passages 95 and 96 are connected to each other to form a single flow passage 97. In addition, an opening of the lower first flow passage 95 and an opening of the lower second flow passage 96 are located in the same position in the axial direction, and these flow passages 95 and 96 are connected to each other to form a single flow passage 97. For preventing fluids flowing through the upper and lower flow passages 97, 97, respectively, from leaking to the outside, a plurality of seal rings 93 are provided between the case body 91 and the shaft body 92."}
{"text": "The field of integral images has developed fast during the last years. Integral image devices are today used for creating eye-catching visual effects for many different purposes, e.g. as security markings, tamper indications or simply as esthetic images. Usually, the integral image device is intended to be provided as a label or as an integrated part in another device. Many different optical effects have been discovered and used and often different optical effects are combined to give a certain requested visual appearance. At the same time, it is a general request to keep the integral image device as thin as possible, in order to facilitate the integration into e.g. other sheet materials.\nThe typical approach for an integral image device is to provide an array of small focusing element. The focusing element may be different kinds of lenses, apertures or reflectors. An image fragment plane is provided with image fragment structures. The image fragment plane is provided relative to the array of focusing elements such that when the device is viewed from different angles, different parts of the image fragments structures are enlarged by the focusing elements and together form an integral image. Depending on the design of the image fragment structures, the integral image can change in different ways when the viewing conditions, e.g. viewing angles, are changed.\nA typical realization of an integral image device is a thin polymer foil, where focusing elements and image fragment structures are created in different planes.\nIn the published European patent application EP2335943, the use of a reflective layer within the foil of an integral image device is disclosed. The image fragment structures are provided directly below the array of micro-lenses and the reflective layer faces the image fragment structures and the micro-lenses from below. The image fragment structures are seen by a viewer as an image reflected in the reflective layer and refracted in the micro-lenses. The optical length between the micro-lenses and the image fragment structures should typically be in the vicinity of a focal length of the micro-lenses in order to give an integral image. In such a design, the optical path is about double the distance between the bottom part of the micro-lenses and the reflective layer. In this way, the total thickness of the integral image device can be reduced.\nHowever, since more elaborate optical designs typically requires additional layers and thereby increase the thickness of the integral image device, there is still a need to further reduce the device thickness."}
{"text": "Mammography has long been used to screen for breast cancer and other abnormalities and for diagnostics. Traditionally, mammograms were formed on X-ray film, but more recently flat panel digital imagers have been introduced that acquire a mammogram in digital form and thereby facilitate analysis and storage and provide other benefits as well. Further, X-ray tomosynthesis of the breast has been proposed recently, as discussed in the earlier-filed applications identified above, and clinical testing has been carried out. The assignee of this patent specification, Hologic, Inc., has demonstrated at trade shows in this country a fused, multimode mammography/tomosynthesis system that takes either or both types of images, and either while the breast remains immobilized or in different compressions of the breast. Dedicated breast tomosynthesis systems also have been proposed.\nTomosynthesis as used in the systems and methods disclosed in this patent specification typically involves acquiring a plurality of tomosynthesis projection images Tp at respective angles relative to the breast, and reconstructing therefrom a plurality of tomosynthesis reconstructed images Tr representative of breast slices that have selective thicknesses. Proper display techniques are desirable to make the presentation of Tp and/or Tr images (collectively referred to here as T images) more effective and efficient for review by health professionals. When tomosynthesis projection images Tp are acquired along with conventional 2D mammograms Mp, improved display methods are desirable that facilitate the display of both T and Mp images. Effective display approaches also are desirable when tomosynthesis images Tp and/or Tr that are acquired at one time need to be compared to mammograms Mp and/or to tomosynthesis images Tp and/or Tr acquired at a different time. Effective displays also are desirable when only Tr and/or Tp images are being displayed. Another display issue relates to Computer Aided Detection (CAD) methods that use computer analysis of images to identify locations and possibly other characteristics of suspected abnormalities. CAD marks currently are placed on or are otherwise associated with mammogram images Mp, but it may be useful to place them at the appropriate location on Tr and/or Tp images o to otherwise associate them with Tr/Tp images. Conversely, it may be desirable to obtain CAD marks by processing Tp and/or Tr images, and place them at appropriate locations on Mp images. Here the notation Mp refers to a conventional mammogram, which is a two-dimensional projection image of a breast; the term Mp encompasses both a digital image as acquired by a flat panel detector or another imaging device and the image after conventional processing to prepare it for display to a health professional or for storage, e.g. in the PACS system of a hospital or another institution. Tp refers to an image that is similarly two-dimensional but is taken at a respective tomosynthesis angle between the breast and the origin of the imaging X-rays (typically the focal spot of an X-ray tube), and also encompasses the image as acquired as well as the image after being processed for display or for some other use. Tr refers to an image that is reconstructed from images Tp, for example in the manner described in said earlier-filed patent applications, and represents a slice of the breast as it would appear in a projection X-ray image of that slice at any desired angle, not only at an angle used for Tp or Mp images. The terms Tp, Tr and Mp also encompasses information, in whatever form, that is sufficient to describe such an image for display, further processing, or storage. The images Mp, Tp and Tr typically are in digital form before being displayed, and are defined by information identifying properties of each pixel in a two-dimensional array of pixels. The pixel values typically relate to respective measured or estimated or computed responses to X-rays of corresponding volumes in the breast (voxels or columns of tissue).\nYet another issue concerns the large storage requirements of tomosynthesis images Tp and/or Tr. Because the reconstructed datasets for Tr images are large, it may be better in some circumstances to store unreconstructed projections images Tp, which require less storage. Transmission times to the storage device, and from the storage device to the display workstation, can thus be reduced. The Tp images in this case can be reconstructed to Tr images just prior to viewing that requires a display of Tr images. Further, it may be desirable that images viewed on a workstation are the same or at least comparable to images viewed on a different workstation, or the same or at least comparable to previously viewed images of the same dataset, even if the software and/or hardware of the workstation or acquisition system or intermediate storage or processing systems, have changed.\nYet another issue concerns the processing time required to reconstruct tomosynthesis images Tr. Because of relatively long reconstruction times, one possible approach is to perform reconstructions at an acquisition console, and send the already-reconstructed images to display workstations. This can allow a greater reading throughput if there are several acquisition systems that are all pushing images to one or more display workstations. The system can be designed so that it can handle M acquisition consoles sending their images to all N display workstations.\nIn one non-limiting example disclosed in this patent specification, acquisition and display of x-ray images starts with acquiring x-ray mammography image data representative of projection mammography images Mp of patients' breasts and x-ray tomosynthesis image data representative of projection images Tp taken at different angles of at least a source of imaging x-rays relative to the patients' breasts (e.g., different angles of the focal spot in an X-ray tube relative an immobilized breast). This acquisition can be performed by a single unit, using a single X-ray tube and a single flat panel digital imager or some other imaging device, configured to selectively acquire one or both of the mammography and tomosynthesis image data, in the same compression of a patient's breast or in different compressions, and at the same imaging session or at different times. The disclosed system and method use at least a subset of the acquired Tp images to form reconstructed tomosynthesis images Tr representative of slices of the breasts that have selected orientations and thicknesses. The system and method display at least a selected subcombination of the Mp, Tr and Tp images, preferably for concurrent viewing or at least for viewing in a single session, and preferably but not necessarily while showing, at or near the displayed images, respective labeling symbols identifying them as Mp, Tr or Tp images and possible other information that facilitates detection/diagnosis of abnormalities, such as information indicative of the position and orientation of the slices represented by Tr images, the thicknesses of such slices, etc. The information can be in alphanumeric or non-numeric form such as in the form of graphics and/or icons.\nThe method and system can further generate or otherwise obtain computer aided detection (CAD) marks for suspected abnormalities in said Mp images, and can display said marks at corresponding locations on or in association with Tr and/or Tp images that are related, e.g. by orientation or otherwise, with respective Mp images. In addition to location information, the CAD marks can provide information regarding, for example, the type of suspected abnormality and/or a confidence level that a mark points to an actual abnormality. CAD marks that are initially generated from or are otherwise related to some of the Tr, Tp or Mp images can be displayed at or in association with images from which they were not generated or with which they were not initially associated, at corresponding or at least related locations. Tp images can be stored together with version information indicative of at least one of an acquisition configuration used to acquire them and a reconstruction configuration used to reconstruct Tr images from said Tp images, to thereby enable later reconstruction of Tr images that match those reconstructed at an earlier time. Alternatively, or in addition, Tp images can be stored together with version information related to when they were acquired and can be later reconstructed into Tr images using a reconstruction configuration that matches the version information. A reconstruction configuration can be provided that has at least two different versions of reconstruction software, so that Tr images can be reconstructed using a version of the reconstruction software that matches the version information of the Tp images or earlier Tr images. Tr images can be reconstructed from only a subset of the acquired Tp images, which in an extreme case means reconstruction from a single Tp image to yield a Tr image that is equivalent of the Tp image. Tr images representative of at least two breast slices that differ in thickness can be formed, for example using MIP (Maximum Intensity Projection) methods or a summing method that may or may not use different weighting of the summed pixel data. The display can be toggled between Tr images representative of breast slices having different thicknesses, wherein the slices may or may not overlap in space. Through computer-processing, the volume of a lesion can be computed and displayed from information contained in the Mp, Tr and/or Tp images. The display can show concurrently, for example, Tr images reconstructed from a current acquisition of Tp images and at least one Mp image obtained from a previous acquisition involving a different breast compression. The concurrent display can be on the same or different display monitors, and can include at least Mp and/or Tr images, or at least Mp and/or Tp, images, or at least Tr and/or Tp images, or all three types of images, and can instead or additionally include 3D images formed from some or all of the acquired X-ray data, image data and/or from Mp, Tr and/or Tp images. Information indicative of status of loading Tr images for display can be shown as a part of the display. Different images can be displayed at different pixel sizes or fields of view or, alternatively, they can be selectively equalized by pixel size or field of view by selected types of interpolation or extrapolation, for example by up-converting to a smaller pixel size and thus a higher converted pixel count or by down-converting to a larger pixel size and thus a lower pixel count.\nThe display of Mp and T images such as Tr images can include displaying non-numeric indications of various properties of the displayed images. Such non-numeric indications can include indications of respective levels, spacing and slice thickness for displayed Tr images relative to a compressed breast imaged in an Mp image or relative to some other frame of reference, for example in the form of cross-lines on a bar related to the breast as displayed as an Mp image, wherein the height of the bar may relate to the thickness of the compressed breast, and/or non-numeric indications of respective thicknesses of breast slices represented by displayed Tr images, for example in the form of cross-bars of respective thickness on a bar related to Mp images, and/or non-numeric indications of the inclination angle of the slice or slices represented by one or more Tr images relative to a selected frame of reference such as the compressed breast that was imaged to generate the data from which the Tr images were reconstructed, and/or non-numeric indications of other parameters. Instead, or in addition to such non-numeric indications, numerical indications can be provided and displayed of the position of a slice image Tr relative to, e.g., a breast imaged in an image Mp, the thickness of the slice represented in a Tr image, and/or the orientation of that slice. Mp and T images such as Tr images can be shown overlaid on each other, and toggling can be allowed to switch between the image that is visible at the time and an image that is not, i.e. toggling between the 2D Mp image and the 3D Tr image or images. Similar toggling is available between different types of T images, of between different T images of the same type. In addition, other image display effects can be provided, such as, without limitation, fade-in/fade-out and blending two or more images at respective weightings, image overlays and masking, or other effeces, as commonly used in post-production of television images and in known image processing software such as Photoshop from Adobe.\nT images such as Tr and/or Tp images can be displayed in ciné mode, with selective control over the speed of changing from one image to another, the order of images for display relative to an order in which they were reconstructed or acquired, the selection of the first and last images in the cine sequence, and/or other parameters. For example, at least two sets of Tr images, each set reconstructed from different acquisitions of Tp images, can be shown concurrently and scrolled through and/or displayed in cine mode in synchronism. For example, the two sets of Tr images can be synchronized such that the first and last slice images of each set can appear on the display at the same time, which may be implemented in different ways, such as by making the number of images the same (but possibly representing slices of one thickness in one set and slices of a different thickness in the other set), or by scrolling through images of one set faster than through images of the other set, or in some other way. Making the number of Tr images of one set match that of the other set may involve interpolating images of new slices to either reduce the number of images/slices in one set or increase it. The images in that case in one set of Tr images can be made to represent slices that are all the same thickness or slices that include some that are thicker or thinner than other slices.\nTr images can be displayed in scroll or cine mode in different orders, such as by starting with the image representing the bottom of the breast (the breast side resting on a breast platform during imaging) toward to image representing the top of the breast (the breast side compressed by a compression paddle), or in reverse order, or by starting with a Tr image representing a selected intermediate slice and proceeding toward the top or the bottom of the breast, or in some other desired order. The initial image that is displayed can be the bottom image, the top image, or a selected intermediate image. In a scroll or cine mode, the display can show every image of a Tr set for a breast, or a selected subset, such as every other image or some other subset of all Tr images representing a breast. The display system can have a default mode for a new user in which it starts the scroll or cine display with the bottom image, but with provisions for this default to be changed to another display protocol for that user in which the starting image is another one of the Tr images. A selection of initial or default display modes can be provided relating to the order, speed, slice thickness and/or other parameters of display of images, and user selection among those modes can be allowed. Similar procedures are available in the case when the orientation of the Tr slices is not the same as that of the Mp or Tp slices.\nThus, information regarding image data acquisition, storage, reconstruction and/or other parameters can be selectively displayed. For example, Tr images can be identified on the display as such, thus providing an image type indication display, and additional identification can be displayed in association with a displayed image that identifies the position in the breast of the slice represented by the displayed Tr image, the thickness of the slice, the orientation of the slice, and/or some other property associated with the displayed Tr image. Corresponding or at least similar display of information can be associated with different types of displayed images, such as Tp and Mp images and displayed in association with the display of the images. In the alternative, images such as Tr, Tp, and/or Mp can be displayed in different combinations and/or sub-combinations that may include the same type of images displayed concurrently or toggled, or different types of images displayed concurrently or toggled, without providing some or any of the identifications discussed above. Tr images can be printed in an N×M format (where N and M are positive integers), and printing of any images displayed on one or more monitors in WISIWIG format can be provided.\nCompression of Mp, Tp, and/or Tr images and/or of other image data can be selectively carried out prior to storing or archiving the images. The compression can be lossless, or it can be lossy to a selected degree. Reconstruction of Tr images can be selectively carried out from compressed Tp images, preferable after suitable decompression. Window/level controls can be provided for at least selected ones of the displayed Tr images, and the controls can be set by the user, or automatically, to control the window width and/or the window level of only one, or only selected ones, or all of the displayed images. Image regions can be magnified for display, and the window/level controls automatically applied to the magnified regions. The Tr, or the Tp, or both the Tr and Tp images, can be stored in PACS storage, and can be associated with related Mp images and/or with selected CAD information. The Tp images can be acquired by using coarser binning, e.g. in a direction of relative motion between the source of imaging x-rays and a breast during image acquisition, or in both directions. Alternatively, such binning can be done after the Tp images are acquired, to thereby reduce storage and further processing requirements. The Mp, Tr, and/or Tp images that are concurrently displayed can be based on image data acquired from the same breast of a patient while the breast remains immobilized under compression that can remain the same or change between the acquisition of Mp images and Tp images. Alternatively, the Mp and Tp images can come from different acquisitions at different times or different breast compressions. Image data for Tp images acquired at two or more acquisition units can be supplied to and reconstructed into Tr images at a single reconstruction unit, from which one or more data display units can acquire Tr, Tp, and/or Mp images for display, or image data for Tp images can be stored as such and only reconstructed into Tr images immediately prior to display thereof. Images with difference characteristics such as pixel size, brightness, gamma curves, etc. can be processed to make selected ones of their characteristics sufficiently similar to facilitate comparison.\nAn additional or alternative display approach uses the Tp and/or Tr images in stereoscopic display. For example, when any two Tp images taken at different angles to the breast are displayed concurrently and viewed such that each is seen by a different eye of the observer, depth information is visualized. Similarly, when any two Tr images are reconstructed such that their image planes are at an angle to each other, depth information can also be perceived.\nImages of different types and from different sources can be displayed in desirable size and resolution. For example, an image can be displayed in (1) Fit To View Port mode, in which the size of the displayed image size is maximized such that the entire imaged breast tissue is visible, (2) True Size mode, in which a display pixel on the screen corresponds to a pixel of the image, or (3) Right Size mode, in which the size of a displayed image is adjusted so that it matches that of another image that is concurrently displayed or with which the displayed image is or can be toggled. For example, if two images of the same breast are taken and are not the same size or do not have the same special resolution, provisions are made to selectively zoom in or zoom out one of them, or zoom both, such that they appear to be the same size on the screen when they are concurrently displayed or the user toggles between them, to facilitate comparison or to otherwise facilitate detection/diagnosis. Known interpolation/extrapolation and weighting techniques can be used in such re-sizing, and known image processing technology can be used to make other characteristics of the displayed images similar in a way that facilitates detection/diagnosis.\nSelected hanging protocols are provided that are unique to the different types of images with which the disclosed system deals. As one example, the hanging protocols for 2D images (e.g. Mp images) and 3D images (e.g. Tr images) are linked so that when one type of image is displayed for a given breast the other type is displayed as well. For example, when the Mp image of a breast is displayed, a tile of the Tr images and/or of the Tp images is automatically displayed at the same time, with a desired hanging protocol that may involve scrolling or cine mode presentation, or may require user input so select a particular subset of the Tr and/or Tp images or a particular individual Tr/Tp image. Thus, a combined hanging protocol set can be provided for 2D and 3D images that are concurrently displayed or toggled such that only one type is displayed at one time. In addition, the combined hanging protocol can include provisions for linked display of CAD information associated with one or both of the 2D and 3D images involved in the hanging protocol. Alternatively, the hanging protocols for 2D images are made different from those for 3D images.\nWhen CAD information is available that is associated with any of the images that can be displayed, linking can be provided such that CAD information derived from one of the types of displayed images can be automatically associated with corresponding locations on another displayed image. For example, if CAD information has been derived on the basis of an Mp image and the Mp image is displayed at the same time as, or is toggled with, Tr images, provisions are made to selectively display CAD information associated with the appropriate locations on Tr images. Alternatively, if CAD information is derived on the basis of Tr and/or Tp images, it can be automatically selectively displayed in association with an Mp image that is displayed at the same time or is toggled with the Tr/Tp images. In one example, an Mp image with CAD marks thereon remain displayed on the screen while Tr images of the same breast are scrolled or shown in cine mode on the screen, to facilitate identification and/or assessment of suspected abnormalities in the Tr images.\nWhen Tr images are displayed, provisions are made to selectively adjust the thickness of the slice represented by any displayed Tr image. For example, if it is desired to display a Tr image of a slice (slab) of breast tissue that is 1 cm thick but the available Tr images represent 1 mm thick slices, 10 of those Tr images can be combined into a single new Tr image that represents the 1 cm slice, using for example a known MIP (maximum intensity projection) technique.\nFor storage, transmission to remote locations, and/or other purposes, the images can be formatted consistent with DICOM standards. For example, each raw or displayed projection image set or reconstructed slice image set for a single view is stored as a single Secondary Capture image instance according to DICOM. The image pixel data can be encoded in a selected compressed format (CODEC) that includes all projection or slice images"}
{"text": "The subject matter disclosed herein relates generally to turbines and, more specifically, to integrating sections of a turbine system.\nA turbine system may include an exhaust diffuser section coupled to a turbine section downstream of the turbine section. Such a turbine system may be either a gas turbine system or a steam turbine system. Specifically, a gas turbine system combusts a mixture of fuel and air to generate hot combustion gases, which in turn drive one or more turbines. In particular, the hot combustion gases force turbine blades to rotate, thereby driving a shaft to rotate one or more loads, e.g., electrical generators, and so forth. The exhaust diffuser section receives the exhaust from the turbines, and gradually reduces the pressure and velocity of the exhaust. Certain turbine systems include a turbine section and a diffuser section that are independently designed for optimal performance. Unfortunately, when such systems are integrated, the combined turbine section and diffuser section may not function optimally."}
{"text": "1. Field of the Invention\nThe present invention relates to integrated circuit testing. In particular, the invention relates to use of a Joint Test Action Group (JTAG) interface in integrated circuit testing.\n2. Description of the Prior Art\nThe use of various functional components (blocks) and component libraries to create complex integrated circuits is becoming increasingly common. These complex integrated circuits are difficult to test. Traditionally, these integrated circuits incorporate a JTAG test port as a mechanism for emulating and debugging an integrated circuit under test. The JTAG test port is used to download debugging or other testing software from an external device, e.g., a JTAG debugger device, to the functional blocks located on the integrated circuit under test. The JTAG test port is also used to upload the test results from the functional blocks to the JTAG debugger device. A JTAG (test port) standard has been adopted by the Institute of Electrical and Electronics Engineers and is now defined as IEEE Standard 11491, IEEE Standard Test Access Port and Boundary-Scan Architecture, which is incorporated herein by reference. See also C. M. Mannales, and R. E. Tullos, xe2x80x9cThe Test Access Port and Boundary-Scan Architecture,xe2x80x9d (IEEE Computer Society Prep, 1990) which is also incorporated herein by reference.\nThe JTAG test port provides control and observability of the functional blocks (test points) of the integrated circuits through scan access. Essentially, a large shift register referred to as a data chain is provided which allows data to be written to or read from desired test points inside the integrated circuit. These data chains are often controlled by a JTAG controller located external to the integrated circuit. The JTAG controller provides the necessary data and clock signals and an interface to a user or an application for automated testing of the functional blocks within an integrated circuit.\nAs shown in FIG. 1 each of the functional blocks (core devices e.g. RAM 101, ROM 103, memory 105) located on a typical integrated circuit 100 requires a particular JTAG test port (interface 107 for connecting to debugger devices 110). Each JTAG test port may be a software or a hardware interface. In either case, it is designed specifically for testing the functions of the corresponding functional block. Furthermore, each of the JTAG test interfaces comprises a plurality of hardware pins (up to five hardware pins for each interface).\nEach hardware pin takes up physical space on the integrated circuit and the circuit board where the integrated circuit is located. Thus, the cost associated with the manufacturing of integrated circuits and circuit boards increases with the increased number of hardware pins. Also, additional hardware pins increase the static on an integrated circuit 100. This static causes interference in the functionality of the integrated circuit. Therefore, it is desired to minimize the number of hardware pins; however, current technology dictates that there must be a particular JTAG interface for each functional block located on an integrated circuit (under test). Each of these JTAG interfaces further requires a plurality of hardware pins.\nOne known solution for reducing the number of hardware pins is to utilize a common debugging platform that allows a plurality of functional blocks to share a common JTAG software or hardware interface. One such known debugging platform is Multi-ICE manufactured by Advanced RISC Machine, 985 University Avenue, Suite 5, Los Gatos Calif. 95030. The ARM Multi-ICE debugging platform can reduce the number of hardware pins, but the use of such a debugging platform limits the integrated circuit testing to the test software written specifically for the common debugging platform. This prevents using off-the-shelf JTAG debugger devices that a user may want or already have. It also creates the need for each vendor to develop different testing software in accordance with each particular debugging platform.\nThe present invention is directed to a JTAG port-sharing device that reduces the required number of hardware pins on an integrated circuit without limiting the integrated circuit testing to a particular debugging platform. The JTAG port-sharing device eliminates the need to have individual hardware pins for each functional block that requires a JTAG test interface and allows one set of hardware pins to be shared by a plurality of functional blocks. The JTAG port-sharing device is not limited by a particular debugging platform and allows the use of pre-owned and/or off-the-shelf testing software and/or JTAG debugger devices thereby leading to shorter development cycles and lower development costs.\nIn one embodiment, the present invention is a port-sharing device for testing an integrated circuit having a plurality of functional blocks, wherein the port-sharing device comprises (a) an on-chip interface port configured to be connected to a pin on the integrated circuit; and (b) at least two debugger ports, each configured to be connected to at least one debugger device to enable the at least one debugger device to test the plurality of functional blocks via the pin.\nIn another embodiment, the present invention is a method for testing an integrated circuit having a plurality of functional blocks, comprising the steps of (a) configuring an on-chip interface port of a port-sharing device to a pin on the integrated circuit; (b) configuring at least two debugger ports of the port-sharing device to be connected to at least one debugger device; and (c) testing the plurality of function blocks in the integrated circuit with the at least one debugger device via the port-sharing device and via the pin."}
{"text": "This invention relates to a wet shaving razor comprising a cartridge that includes a shaving blade with a cutting edge which is moved across the surface of the skin being shaved by means of an adjoining handle. Conventional safety razors have a blade unit connected to a handle for a pivotal movement about pivotal axis which is substantially parallel to the blade or the blade edge. For example, U.S. Pat. Nos. 7,197,825 and 5,787,586 disclose such a razor having a blade unit capable of a pivotal movement about a pivot axis substantially parallel to the blade(s). The pivotal movement about the single axis provides some degree of conformance with the skin allowing the blade unit to follow the skin contours of a user during shaving. Such safety razors have been successfully marketed for many years. However, the blade unit can fail to remain flat and often disengages from the skin during shaving due to the blade unit's limited ability to pivot about the single axis combined with the dexterity required to control and maneuver the razor handle. The combination of these deficiencies can affect the glide and overall comfort during shaving.\nThere have been various proposals for mounting a cartridge on a handle to enable movement of the cartridge during shaving with the aim of maintaining conformity of the skin contacting parts with the skin surface during shaving. For example, many razors currently marketed have cartridges which are pivotable about longitudinal axes extending parallel to the cutting edges of the elongate blades incorporated in the cartridges. There is an increasing need to provide a shaving consumer with a closer, more effective shave. Applicant has attempted to provide this in its commercially available Fusion® razor which incorporates a spring in its following system to bring about a reduced cartridge to skin angle, which has been found to lead to a better shave. Similarly, others have attempted to manipulate the biasing mechanisms of their commercial razors. For instance, US Patent Publication 2005/0241162 A1 discloses a biasing assembly for a wet shave razor wherein the assembly includes 1) an abutment surface defined by a cartridge and located on the underside of the cartridge and 2) a biasing member extending outwardly from the handle and having an end which when the cartridge is coupled to the handle is in sliding engagement between the neutral and fully-rotated positions. The biasing member exerts a variable torque against the abutment surface. The reference, however, focuses primarily on a low spring force to prevent the cartridge from lifting off of the skin and does not focus on the effect that the biasing member has on maintaining the cartridge flat relative to the skin during shaving strokes and corresponding shaving closeness.\nIn addition, current shaving razors found on the market typically include handle configurations that are variations of an ‘L’ shape where the longitudinal axis 30 of the handle 14 is offset from the razor cartridge 12 such that it intersects the cutting plane 122 behind the cartridge 12 as shown in FIG. 1. This configuration has the effect of pushing the razor cartridge 12 through the shaving stroke which can make it difficult to maneuver and can require a steady hand to steer the razor cartridge 12. In addition, the shaving razors have an axis of roll 36 (interchangeably referred to hereinafter as axis of roll 36 and handle roll axis 36) that extends between the free end of the handle 14 and a point on the cartridge where the forces are balanced. The axis of roll 36 is the line about which the razors spin in the direction shown in FIG. 1 when in a user's hand. For the L-shape configuration shown, this arrangement has a shortcoming. Since the handle longitudinal axis 30 extends above the axis of roll 36, instability is introduced during shaving, similar to a top heavy scenario that a user must compensate for when handling the razor. Hence, additional effort is required by the user to maintain stability of the razor during shaving.\nIn pursuit of an improved shaving product, there is a need for a shaving razor that can maintain the blade unit of a razor cartridge flat against the skin throughout a shaving stroke. Particularly there is a need for a shaving razor having a biasing member producing a progressively increasing return torque on a cartridge forcing the cartridge into contact with the skin throughout the shaving stroke. In addition, there is a need for a handle geometry that provides the user with improved control while shaving."}
{"text": "1. Field\nThe present disclosure generally relates to a light-emitting diode (LED) driver circuit for use with LED bulbs, and more particularly, to and LED driver circuit with an improved power factor.\n2. Description of Related Art\nDespite the many benefits of LED bulbs, there are some challenges that have prevented LED bulbs from widely replacing incandescent and fluorescent bulbs in residential application. For example, electrically, LED bulbs operate differently than incandescent and fluorescent bulbs. LED bulbs are current controlled devices, meaning that the light output is control by changes in current as opposed to incandescent and fluorescent bulbs that are voltage controlled.\nThe difference in control requires that LED bulbs have special driver circuits that convert the standard AC voltage supplied in residential outlets to a current suitable for driving LEDs. These driver circuits, however, typically result in an LED bulb that interacts with the electrical grid very differently than incandescent bulbs.\nPower factor is one significant parameter where LED bulbs differ from incandescent bulbs. Power factor is the ratio of real power flowing to a load to the apparent power. A load with a power factor of 1 means that the load is using all power being delivered to the load. Typically, purely resistive loads have a power factor of 1. A power factor of less than 1 indicates that there is energy storage in the load that may return power to the power supply out of phase with the power supply. The lower the power factor, the more wasted power.\nLED bulb driver circuits typically have storage elements (e.g., capacitors) that may cause a lower power factor for the LED bulb as compared to an incandescent bulb. This results in an LED bulb that may put more strain on the power supply (i.e., the electrical grid) than is necessary.\nLED bulb driver circuits may be modified with additional components or special circuits to improve the power factor. However, these modifications increase the volume occupied by the driver circuit. In space limited LED bulbs, it may be difficult to fit these additional components or special circuits. Additionally, the modifications may also make it more difficult for the LED bulb to work with common residential light dimmers."}
{"text": "The main wing of an aircraft has so far been produced using an aluminum alloy or a titanium alloy, but recently, has used a carbon fiber-reinforced plastic (CFRP) comprising a combination of carbon fibers and resin.\nIn the aircraft, damage to the wing poses a danger directly leading to a serious accident. Thus, it is very important to perform operations for maintenance and checkup, thereby finding a defect.\nDefects occurring in the wing of the aircraft include a flaw, foreign matter, a peel, and so on. As a nondestructive examination for detecting these defects, an ultrasonic flaw detection test or inspection is adopted.\nA general technique for the ultrasonic flaw detection test will be explained first of all.\nIf a flaw detection range A is broad, as shown in FIG. 10, a flaw detection path P is established or set. Along this flaw detection path P, an ultrasonic probe is moved to acquire a flaw detection waveform signal (response waveform signal) by an ultrasonic wave.\nIf, at a certain point on the flaw detection path P, this site is a sound site, a flaw detection waveform signal, which comprises a surface echo Es and a bottom echo Eb as shown in FIG. 11, is obtained. If, at a certain point on the flaw detection path P, this site is a defective site, a flaw detection waveform signal, which comprises the surface echo Es, the bottom echo Eb, and a defect echo Ed as shown in FIG. 12, is obtained.\nIf the flaw detection range A is narrow, it is possible to detect a defect (detect the defect echo Ed) by visually observing the waveform of the flaw detection waveform signal throughout the flaw detection path P. If the flaw detection range A is wide, on the other hand, the visual confirmation of the waveform of the flaw detection waveform signal over the entire flaw detection path P takes a huge time.\nUnder these circumstances, the flaw detection waveform signal is subjected to signal transformation or signal conversion to convert this signal into a flaw detection image signal representing a flaw detection image in which the luminance or brightness value of the defective site changes relative to the brightness value of the sound site. By visually confirming the flaw detection image based on this flaw detection image signal to determine whether there is a defect or not, a flaw detection examination or test is conducted in a short time.\nIf the presence of the defect is observed, a detailed analysis can be made by confirming the flaw detection waveform which is the basis for the flaw detection image.\nAs techniques for signal transformation for converting the flaw detection waveform signal into the flaw detection image signal, TOF image signal transformation and AMP image signal transformation are known.\nTo carry out signal transformation, a gate G with a predetermined range and at a predetermined height is set at a position between the surface echo Es and the bottom echo Eb in each of FIG. 13 showing the flaw detection waveform signal on the sound site and FIG. 14 showing the flaw detection waveform signal on the defective site.\nFor conversion into a TOF image, the position where the set gate G and the defect echo Ed intersect is converted into a brightness value. For conversion into an AMP image, the height of the defect echo Ed intersecting the set gate G is converted into a brightness value.\nIf a plurality of the defect echoes are present, there are versions, for example, in which the first echo that intersects the gate G is used as the defect echo Ed, or the highest echo of the echoes intersecting the gate G is used as the defect echo Ed.\nFIG. 15 shows an example of the flaw detection image of the sound site. FIG. 16 shows an example of the flaw detection image (TOF image) of the defective site. FIG. 17 shows an example of the flaw detection image (AMP image) of the defective site.\nAs shown in these drawings, the flaw detection image shows the brightness of the defective site changing relative to the brightness of the sound site, and enables the defective site to be inspected visually with rapidity and ease.\nThe waveform shape of the flaw detection waveform signal and the status of occurrence of a noise differ according to the type of a defect or the plate thickness, shape, and material of the flaw detection site. Thus, a single flaw detection image may be insufficient to determine accurately whether there is a defect or not.\nThus, a plurality of gates different in position or height are set, in accordance with the type of the defect or the plate thickness, shape, and material of the flaw detection site, to generate a plurality of flaw detection image signals, and a plurality of flaw detection images by the plurality of flaw detection image signals are visually confirmed, whereby the presence or absence of the defect is determined overall.\nThat is, one flaw detection waveform signal obtained when the probe is moved along the flaw detection path P is subjected to signal transformation by a plurality of different signal transformation techniques to obtain a plurality of flaw detection image signals, and flaw detection images by these plural flaw detection image signals are used to judge globally whether the defect is present or not.\nConversion of a flaw detection waveform signal into a TOF flaw detection image signal, and conversion of a flaw detection waveform signal into an AMP flaw detection image signal are different techniques for signal transformation. Moreover, the conversion of a flaw detection waveform signal into a TOF flaw detection image signal, with the gate G differing, is also a different technique for signal transformation. Furthermore, the conversion of a flaw detection waveform signal into an AMP flaw detection image signal, with the gate G differing, is also a different technique for signal transformation. In these senses, these methods are called “a plurality of different signal transformation techniques”.\nA general example of globally judging the presence or absence of the defect by the flaw detection images by the plurality of flaw detection image signals obtained by the signal transformation of the one flaw detection waveform signal will be described below.\nIf the types of the defect are two types, D1 and D2, the plate thickness of the flaw detection site is 10 mm, the materials of the flaw detection site are of two types, M1 and M2, and the shape of the flaw detection site is of two types, F1 and F2, for example, a flaw detection image I1 by a TOF flaw detection image signal obtained by signal transformation of the flaw detection waveform signal with a first gate being set, a flaw detection image I2 by a TOF flaw detection image signal obtained by signal transformation of the flaw detection waveform signal with a second gate being set, a flaw detection image I3 by an AMP flaw detection image signal obtained by signal transformation of the flaw detection waveform signal with a third gate being set, and a flaw detection image I4 by an AMP flaw detection image signal obtained by signal transformation of the flaw detection waveform signal with a fourth gate being set can be used for defect evaluation.\nIf the allocation of the flaw detection images I1 to I4 is made as in Table 1, the flaw detection images I1 to I4 are visually confirmed in sequence, whereby all the defects can be evaluated.\nTABLE 1Example of allocation for defect evaluationby flaw detection imagesPlate thicknessMaterial ofFlawType ofof flawflaw detectionShape of flawdetectiondefectdetection sitesitedetection siteimageD110M1F1I1D110M1F2I1D110M2F1I2, I3D110M2F2I4D210M1F1I1D210M1F2I2D210M2F1I1D210M2F2I2, I4\nA conventional technique for conducting the ultrasonic flaw detection test of the main wing of an aircraft will be explained.\nWith the test of the main wing of the aircraft, the whole surface of the main wing is ultrasonically tested for flaw detection. However, it takes too much time to analyze all of flaw detection waveforms by a flaw detection waveform signal obtained by this test. Thus, the acquired flaw detection waveform signal is subjected to signal transformation by a plurality of different signal transformation techniques to be converted into a plurality of flaw detection image signals. On a plurality of flaw detection images by the plurality of flaw detection image signals, a defective site is identified visually.\nOn the flaw detection image, the brightness value of the defective site is increased (or decreased) compared with a sound site. By visually confirming a change in brightness between the sound site and the defective site, therefore, the defective site is identified.\nIn this case, the type and arrangement of the flaw detection image to be visually confirmed differ according to the type of the defect to be detected, or the difference of the site (shape, plate thickness, material). Moreover, the display magnification and contrast of the flaw detection image in which the defect is easy to detect are also different. This is because a manner of reflection of an ultrasonic wave differs according to the type of the defect or the plate thickness, shape or material of each site.\nThus, an inspector repeatedly changes the type and arrangement, display range (magnification) and contrast of the flaw detection image displayed, according to the type of the defect and the difference of the site, and detects the defect by so doing.\nThe details of such a conventional inspection technique will be described by reference to FIG. 18 which is a flowchart, and FIGS. 19 to 21 showing flaw detection waveform images displayed on the display screen of an image display device.\nAssume that a plurality of flaw detection image signals created by signal transformation of a flaw detection waveform signal, which was obtained by the ultrasonic flaw detection test of the main wing of an aircraft, by a plurality of different signal transformation techniques have been previously stored on a flaw detection image signal database.\nIn starting visual inspection work (Step S1 in FIG. 18), an operator opens a plurality of flaw detection images I1 to In derived from a plurality of flaw detection image signals (Step S2).\nThen, the type of a defect to be evaluated (for example, the defect D1 of the first type) is determined (Step S3).\nWhat type of defect should be investigated is judged based on knowledge possessed by the operator (knowledge of the structure, shape, plate thickness and material of the wing, knowledge of the characteristics of the ultrasonic flaw detection test, knowledge acquired through inspection experiences accumulated, and so on).\nThe operator selects one or a plurality of the flaw detection images required, according to the type of the defect determined, and displays the selected image or images in an aligned arrangement on a display screen 1 of an image display device as shown in FIG. 19 (Step S4).\nIf the type of the defect is D1, for example, the flaw detection image I1 and the flaw detection image I3 are selected, and both of the flaw detection images I1 and I3 are displayed in an side-by-side arrangement.\nIn this case, the flaw detection images I1 and I3 are each displayed in such a manner as to be superposed on, for example, partitioned regions R1 to R28 separated into 28 (4×7) regions by parting lines indicated by short dashed lines in FIG. 19. That is, the flaw detection images I1 and I3 are each divided into images on the partitioned regions R1 to R28.\nWhich of the flaw detection images should be selected according to the type of the defect, and how the selected images should be arranged, are judged based on the knowledge possessed by the operator (knowledge of the structure, shape, plate thickness and material of the wing, knowledge of the type of the defect, knowledge of the characteristics of the ultrasonic flaw detection test, knowledge acquired through inspection experiences accumulated thus far, and so on).\nThe operator changes the display range of the flaw detection images I1, I3.\nConcretely, the starting position of X and Y coordinates displayed, and the ending position of the X and Y coordinates are changed to specify the partitioned regions to be displayed. Also, the display magnification is optimally changed to display, on an enlarged scale, an image in a predetermined range, for example, on the partitioned region R1 of the flaw detection images I1 and I3 (Step S5), as shown in FIG. 19.\nTo what degree the display magnification should be adjusted to be increased, and the image on which of the partitioned regions should be displayed, according to the type of the defect, are judged based on the knowledge possessed by the operator (knowledge of the structure, shape, plate thickness and material of the wing, knowledge of the type of the defect, knowledge of the characteristics of the ultrasonic flaw detection test, knowledge acquired through inspection experiences accumulated thus far, and so on).\nIn the state of FIG. 20, the image on the partitioned region R1 has merely been enlarged. Hence, the contrast between an image showing the defective site (a black-colored portion) and the sound site is not clear.\nThus, the operator makes a contrast adjustment for the enlarged images of the flaw detection images I1, I3 on the partitioned region R1 (FIG. 20) to enhance the contrast between the image of the sound site and the image of the defective site (Step S6).\nTo what degree the contrast should be adjusted, according to the type of the defect, is judged based on the knowledge possessed by the operator (knowledge of the structure, shape, plate thickness and material of the wing, knowledge of the type of the defect, knowledge of the characteristics of the ultrasonic flaw detection test, knowledge acquired through inspection experiences accumulated thus far, and so on).\nThe operator visually confirms the images (FIG. 21) obtained by enlarging, and adjusting the contrast of, the flaw detection images I1, I3 on the partitioned region R, to evaluate whether the defective site is present or not (Step S7).\nThe operator changes the display range of the flaw detection images I1, I3 from the image enlarged from the image on the partitioned region R1, successively, to the images enlarged from the images on the partitioned regions R2 to R28. When selecting the enlarged image on each of the partitioned regions R2 to R28, the operator performs the optimal changing of the display magnification, the contrast adjustment, and the visual inspection of the defective site, which are shown in Steps S5, S6 and S7, respectively, in accordance with the situation of the selected partitioned region.\nAfter all the images on the partitioned regions R1 to R28 are inspected for the presence or absence of the defect D1 (Step S8), the type of the defect to be evaluated is changed to the defect of a different type (for example, defect D2) (Step S3).\nFor the defect D2, an inspection of the defective site is conducted in the same manner as for the defect D1.\nIn this manner, the type of the defect is sequentially changed, and after the inspection of the defective site is conducted for all types of defects, evaluation is completed (Step S9). Then, analysis work on the flaw detection images is completed (Step S10).\nDepending on the type of the defect, the analysis of the flaw detection image in a specific partitioned region may be skipped, and this analysis may be made in a next partitioned region.\nThe reason is as follows: Knowledge or experience may teach, from the beginning, that depending on a specific image selected and a specific defect selected, defect detection in a specific partitioned region is impossible. In this case, the analysis of the flaw detection image in the specific partitioned region is skipped in order to cut down on time and labor."}
{"text": "Until recently, voice information was carried only over a network dedicated to the voice information, such as a public switched telephone network (PSTN). Presently, however, other technologies have been developed that enable voice to be carried via a variety of networks that are not exclusively dedicated to the transfer of voice information. One of these technologies is Internet telephony, also called voice over internet protocol telephony (VoIP). VoIP contains technologies that enable voice to be carried via a data network of the Internet. In a VoIP system, the analog signals from the voice of a caller are converted into digital signals by a VoIP device before being transmitted through a gateway. The gateway then routes the signals over the PSTN dependent on the setup of the VoIP provider. The digitized, voice signals accordingly travel as discrete packets over a public or private Internet. This saves costs over a conventional PSTN as there is no toll charge when making a call.\nHowever, various problems exist with VoIP systems. One of these problems is determination of the location of a caller for purposes of emergency (medical, police, fire) assistance. To access emergency assistance, a caller dials 911 in the United States or Canada and is connected to a Public Safety Answering Point (PSAP). The PSAP is one or more entities that receives emergency calls and dispatches the appropriate emergency services, if necessary. As the basic emergency phone number is the same throughout the United States and Canada, the call is transmitted to one or more routers in a PSTN network of routers through a PSTN gateway to determine the appropriate PSAP to respond to the emergency call, that is, the PSAP most proximate physically to the caller. The appropriate PSAP is determined by using an Automatic Location Identification (ALI) database that associates telephone numbers with street addresses and local emergency response. Once the address is determined, it is transmitted to a display for the benefit of the emergency response personnel. Emergency services then react to a 911 call using the location of the caller.\nWhen using a conventional non-portable telephone, the location of the telephone is known and registered as a particular location. However, the information used to route a VoIP call may not be sufficient to determine the location of the VoIP caller. Unlike fixed telephones, which each use a dedicated switch to connect to the PSTN network, a VoIP network uses one switch to serve multiple VoIP devices distributed over a broad geographic area. Moreover, VoIP devices are nomadic, that is, VoIP phones can be transported to any location where a network connection is available, connected to the Internet, and calls may be made from the new location. Thus, there is no fixed relationship between the location of a caller using a VoIP device and the phone number used by the VoIP device. Nor do VoIP networks generally provide information regarding the location of callers.\nThus, there is a need for a VoIP device and associated method in which emergency calls are routed to a proper emergency service provider for providing appropriate emergency response to the caller."}
{"text": "The invention provides a novel class of compounds, pharmaceutical compositions comprising such compounds and methods of using such compounds to treat or prevent diseases or disorders associated with abnormal or deregulated kinase activity, particularly diseases or disorders that involve abnormal activation of the PKC, Jak1, Jak2, Jak3, Tyk2, KDR, Flt-3, ROCK, CDK2, CDK4, TANK, Trk, FAK, Abl, Bcr-Abl, cMet, b-RAF, FGFR3, c-kit, PDGF-R, Syk, or Aurora kinases.\nProtein kinases constitute a large family of structurally related enzymes that are responsible for the control of a variety of signal transduction processes within cells (see, e.g., Hardie and Hanks, The Protein Kinase Facts Book, I and II, Academic Press, San Diego, Calif., 1995). Protein kinases are thought to have evolved from a common ancestral gene due to the conservation of their structure and catalytic function. Almost all kinases contain a similar 25-300 amino acid catalytic domain. The kinases may be categorized into families by the substrate that they phosphorylate (e.g., protein-tyrosine, protein-serine/threonine, lipids etc.). Sequence motifs have been identified that generally correspond to each of these families (see e.g., Hanks and Hunter, (1995), FASEB J. 9:576-596; Knighton et al., (1991), Science 253:407-414; Hiles et al., (1992), Cell 70:419-429; Kunz et al., (1993), Cell 73:585-596; Garcia-Bustos et al., (1994), EMBO J. 13:2352-2361).\nThe protein kinase C family is a group of serine/threonine kinases including at least ten related isoenzymes, including alpha, beta 1, beta 2, gamma, delta, epsilon, eta, lambda, iota, theta and zeta. The isoenzymes have been divided into three groups based on their different expression patterns and co-factor requirements. The classical PKC enzymes (cPKC), including alpha, beta 1, beta 2 and gamma isozymes require diacylglycerol (DAG), phosphatidylserine (PS) and calcium for activation. The novel PKC's (nPKC), including delta, epsilon, theta and eta isozymes, require DAG and PS but are calcium independent. The atypical PKC's (aPKC), including zeta, lambda/iota do not require calcium or DAG.\nPKC isoforms have been shown to play key roles in cellular signaling, proliferation, differentiation, migration, survival, and death. In resting cells, PKCs are predominantly localized in the cytosol and are catalytically inactive due to autoinhibition by their pseudosubstrate domain. Upon cell activation, PKC isotype-specific signals trigger translocation from the cytosol to the membrane and induce conformational changes, which displace the pseudosubstrate moiety from the catalytic domain and enable PKC isotypes to phosphorylate specific protein substrates (Biochem. J. 370:361-371, 2003). Most isoforms are ubiquitously expressed, except PKCγ and PKCθ. While PKCγ is exclusively found in the brain, high protein levels of PKCθ are seen predominantly in hematopoietic cells and skeletal muscle. PKCα and PKCθ as well as PKCβ and PKCδ are functionally important for T and B cell signaling, respectively (Nat. Immunol. 5:785-790, 2004. Curr. Opin. Immunol. 16:367-373, 204. Nature. 416:860-865, 2002). PKCθ plays an essential role in T cell activation because it is the only isoform that is selectively translocated to the T cell/antigen-presenting cell contact site immediately after cell-cell interaction (Nature. 385:83-86, 1997). Furthermore, PKCθ is crucial for IL-2 production, a prerequisite for the proliferation of T cells (Eur. J. Immunol. 30:3645-3654, 2000). PKCθ-deficient mice are defective in NF-κB (Cell Mol. Immunol. 3:263-270, 2006), NFAT and AP-1 activation (Nature, 404 (96776), 402-407, 2000. Journal of Immunology 176:6004-6011, 2006) and are resistant to experimental autoimmune encephalomyelitis (J. Immunol. 176:2872-2879, 2006), collagen-induced arthritis (Journal of Immunology 177 (3), 1886-1893, 2006) and asthma (Journal of Immunology 173 (10), 6440-6447, 2004). PKCα in T cells is required for proliferation and IFN-γ production (J. Immunol. 176:6004-6011, 2006). B cells require PKCβ for proper antigen receptor function and PKCδ for the induction of tolerance (Nature. 416:860-865, 2002). Thus, PKC isoforms in T and B cells are considered attractive therapeutic targets for autoimmune diseases and transplantation (Curr. Opin. Investig. Drugs. 7:432-437, 2006).\nFurther, PKCε and PKCγ have been suggested to play a role in nociception and inflammatory pain (J. Pharm. Exp. Ther. Pain 110, 281-289, 2004) and PKCζ has been proposed as an intermediary in the activation of the NF-κB and IL-4/Stat6 pathway (Cell Death Differ. 13: 702, 2006). The NF-κB pathway is important for inflammatory and immune diseases, therefore a PKCζ inhibition may serve to reduce the severity of these type of diseases (Allergol. Int. 55: 245, 2006. J. Biol. Chem. 281: 24124, 2006. Arthritis Rheum. 56: 4074, 2007. J. Interferon Cytokine Res. 27: 622, 2007).\nThe novel compounds of this invention inhibit the activity of one or more protein kinases and are, therefore, expected to be useful in the treatment of kinase-mediated diseases."}
{"text": "1. Area of Invention\nThe invention is a device, capable of self-stowing, to improve the aerodynamics and/or reduce base drag at the aft end of bluff bodies, such as box shaped heavy trailers towed by powered tractor trucks and similarly shaped cargo and passenger vehicles, that does not impede trailer docking and is capable of incorporating air blowing techniques.\n2. Essential Perameters\nThe purpose of the invention is to improve of the aerodynamics at the aft end of a heavy trailer or other vehicle having a bluff body to reduce aerodynamic drag. The reduction of drag lessens the energy required from, and therefore the fuel consumption, of the powered vehicle or towing vehicle. A device producing this effect is economically beneficial to the operator and will lessen national fuel consumption and emissions.\nThe design of the vast majority of existing and newly constructed heavy trailers incorporates a pair of outwardly hinged doors that cover the entire aft end of the trailer and are capable of swinging through an arc of 270 degrees from the closed position to a position parallel to the outside of the sides of the trailer. This feature is essential to efficient loading and docking but presents a major obstacle to the application of aft end aerodynamic devices.\nAerodynamic devices that need to be removed or partially disassembled to enable door function are cumbersome, costly, subject to loss and can be accidentally dislodged; they require temporary storage and entail lost time and added effort to the operator, to the economic disadvantage of the freight carrier.\nBecause heavy trailers frequently wait in the yard between transits, and because the cost of fuel saved by the device is a finite number, there exists an economic threshold to the cost to manufacture and install the device, precluding complex devices with multiple components.\nTo become widely accepted by the trucking industry an aft end aerodynamic device must be functional, low in cost, easily retrofitted, permanently attached, allow normal door function, and require little or no operator intervention in its use.\n3. Technical Considerations\nExtensive studies have been performed to demonstrate and evaluate the efficacy in drag reduction of various aerodynamic alterations to and of aerodynamic devices for the traditional, rectangular box shape of the heavy trailer and other similarly shaped highway vehicles. Herein, only those modifications relating to the trailing or aft end of the heavy trailer will be addressed.\nIn an SAE (Society of Automotive Engineers) technical paper presented in May 2001 [Advanced Aerodynamic Devices to Improve the Performance, Economics, Handling and Safety of Heavy Vehicles; 2001-01-2072], the author, Robert J, Englar, Principal Research Engineer, Georgia Tech Research Institute, summarizes the results of research conducted for the DOE Office of Heavy Vehicle Technologies. Computer and wind tunnel tests indicated that applying a fixed curved surface to the aft edges of the trailer produced a 7.3% reduction in drag.\nFurthermore, by blowing compressed air tangentially over that surface, the reduction of aerodynamic drag can be significantly enhanced as the blown air entrains air passing along the sides of the vehicle (the so called Coanda Effect); and when the blowing is selectively applied, stability of the trailer can be improved.\nAn AIAA (American Institute of Aeronautics and Astronautics) paper presented in July, 2004 [DOE's Effort to Reduce Truck Aerodynamic Drag—Joint Experiments and Computations Lead to Smart Design; UCRL-CONF-204819], summarizes the results of further DOE sponsored testing of various devices and recites positive drag reduction results of 10-12% from corner rounding and an additional 15% from tangential blowing.\nAccording to a July 2001 article in Trailer-body Builder Magazine [Winds of Change; http://trailer-bodybuilders.com/mag], in which Robert J. Englar was interviewed, The American Trucking Associations determined that each 3% reduction in a vehicle's drag coefficient equates to about 1% in fuel savings. Englar is said to have described the device to be used in testing as “90 degree arcs”, and as “a patented curved surface 6-9″ in radius attached . . . to the four corners of the back end of the trailer” (quotes from the article, not directly from Englar).\nA Georgia Tech Research News release [Jan. 2, 2005, Low-Drag Trucks: Aerodynamic Improvements & Flow Control System Boost Fuel Efficiency in Heavy Trucks; http://gtresearchnews.gatech.edu/newsrelease/truck fuel.htm] further describes the road test devices and their application to vehicle stability as well as drag reduction.\nThe engineering and aerodynamic principles involved in the creation of the cited drag reductions and contributions to trailer stability are fully described in the afore-mentioned engineering papers and need not be recited herein.\nU.S. Pat. Nos. 5,863,090 and 7,192,077 to Englar pertain to racing cars and ground vehicles, respectively, and disclose the theory, application and methods of generation, expulsion and control of air blown over surfaces of a vehicle to create the Coanda Effect with resultant effects on aerodynamic drag and vehicle stability. Devices to enable the production of these effects at the aft of a heavy trailer or other vehicle having a box like shape and bluff aft end is the subject of the invention.\nThe actual rounding of the corners of the trailer is impractical due to considerations of cargo space, loading and the sealing of the rear doors. In light of the above, a need exists for a device to be attached to the aft of the trailer to accomplish the same effect as corner rounding while meeting the essential parameters outlined in [0002-0006] above. The device is to be capable of manufacture in a shape to suit optimal aerodynamic requirements, and to be easily adapted for tangential blowing. [See DOE 21st Century Truck Partnership White Paper, 21CTP-003, pg. 38].\n4. Prior Art\nNumerous aerodynamic devices have been patented that are designed to attach to the rear of the trailer. Because of the narrow scope of the invention, only prior art pertaining to permanently attached devices that simulate the rounding of the corners of the aft end of the trailer, that relate to the blowing of air at the corners or relate to a self-stowing or non-removable concept or system for aft-of-trailer will be considered.\nWhile the means and method of blowing air tangentially over the device of this invention is a part hereof, neither the inventions, per se, of the tangential blowing of air to reduce aerodynamic drag, nor a control system for this purpose, nor a means to generate air for this purpose are a part on the invention, as these are contained in prior art.\nU.S. Pat. No. 6,854,788 to Graham discloses a self-deploying device, which attaches to the rear doors and to the sides of the aft end of the trailer. Flexible flat sheets positioned at the sides of the trailer when the doors are open. As the doors close, the motion of the doors cause the sheets to be forced aft and formed into curved surfaces (airfoils) by means of and with the support of struts acting on them. The device consists of a minimum of 22 individual parts (not counting fasteners).\nWhile these airfoils may serve to reduce aerodynamic drag, they are large and rigid when deployed, making them subject to damage if accidentally contacted. In addition the plurality of parts and interconnections, and multiple fastening points make them costly to manufacture and install. Even if constructed, as suggested, of aluminum, they add considerable weight and the scrap value of this material makes them targets of theft.\nU.S. Pat. No. 5,824,734 to Lechner discloses a design for the aft end of a vehicle causes the top and sides of the vehicle to be angled toward the center at a suggested angle of 10 to 20 degrees. It is unclear whether these surfaces are formed into the body of the vehicle or whether the angled surfaces are rigidly attached extensions. In either case this device is not suitable for application to heavy trailers due to door function and loading requirements, and is not retrofittable.\nLechner further claims devices to blow air at an angle at the outer edge of the rear end of the device, invoking the Coanda Effect, but does not define the configuration or placement of the nozzles to be used to achieve this. It is questionable whether this is a proper application of the Coanda effect.\nU.S. Pat. No. 7,129,077 to Hilleman discloses hollow, rounded devices attached to the vertical edges of the aft of the trailer. These appear at first to be corner rounding devices but are in fact scoops that gather air at the sides and direct it into the space at the aft of the trailer. Although the devices move forward so as to not project behind the aft end of the trailer, there is no provision to allow the doors to be swung flush to the sides of the trailer.\nU.S. Pat. No. 5,280,990 to Rinard discloses a system or combination of individual devices that are to be applied in concert to the Class 8 truck, both to tractor and trailer. The claims include right and left side “arcuate” devices of a deformable, resilient material attached to the aft vertical edges of the trailer to direct air “passing along” the sides of the trailer transversely into the area at the rear of the trailer. These are vanes or deflectors and have no relation to corner rounding.\nThe devices referred to in [0023] are attached to the trailer by articulated support brackets that allow the devices, through the action of opening the doors, to pivot in such a way as to lie between the sides of the trailer and the outer faces of the doors when the doors have been swung 270 degrees from the closed position.\nThe attachment devices referred to in [0024] appear to consist of a minimum of 32 components (not counting fasteners). In the deployed orientation, the vanes and brackets extend aft of the end of the trailer making them susceptible to accidental damage from rear or side impact. Due to these features, the self-stowing aspect of the devices is impractical and requires high cost to manufacture, install and maintain.\nU.S. Pat. App. Pub. No. US2007/0024087 (Skopic) discloses a pair of symmetrical devices consisting of “flexible vanes” attached by cables or “mesh” vertically at both aft corners of the trailer and intended to redirect air “flowing” along the sides of the vehicle into “a volume of low pressure air” aft of the vehicle. The vanes are formed with 2 angles and three faces and approximate a 90 degree short-sided angle with a flatted point.\nThe devices described in [0026] are not intended for corner rounding, but as deflectors. Their ability to self-store between the sides of the trailer and the outer faces of the doors, when opened 270 degrees from the closed position, requires the use of the aforementioned cables or extremely flexible mesh. This makes them subject to damage or dislodgement and reduces their real world practicality.\nU.S. Pat. No. 6,485,087 to Roberge and Boivin discloses a device for attaching flat, hinged, side extensions to the aft vertical and top corners of the trailer and has no relevance to corner rounding. The attaching features of this device use the expedient of replacement of the existing trailer door hinges. Position when deployed is maintained by one or more cables in conjunction with pressure exerted by top flaps, powered by pneumatic springs.\nU.S. Pat. No. 6,257,654 to Boivin, et al. discloses a device for attaching flat, hinged, side extensions to the aft vertical corners of the trailer and has no relevance to corner rounding. The attaching features of this device use the expedient of replacement of the existing trailer door hinges. Position when deployed is maintained through a combination of cables, springs, pneumatic cylinders or the like. A top flap on each door is similarly attached and controlled.\nThe devices described in [0028] and [0029] extend rigidly straight back from the aft of the trailer making them susceptible to accidental damage from rear or side impact. The plurality of parts increases the cost, weight and difficulty of installation.\nThe device described in [0029] differs from that described in [0028] in that the side extensions are shortened at their tops to clear the top of door flaps when stowed. The devices of both these inventions are stowed and deployed by the action of swinging the rear doors 270 degrees from and back to the closed position, requiring the aforementioned springs as a power source, and some sort of restraint.\nU.S. Pat. App. Pub. No. 2007/0089531 (Wood) discloses a device generally referred to a “boattail plates”, and has no relation to corner rounding. The device self stows and deploys by the action of swinging the rear doors 270 degrees from and back to the closed position. A spring hinge, pneumatic cylinder or mechanical brace is employed to hold the device against a stop in its deployed position. Friction occasioned by the trailing edge of the device riding on the outer sidewalls of the trailer during door opening is dealt with through the commonly invoked use of rollers and/or antifriction materials applied to the contact surfaces.\nThe devices described in [0028, 0029 and 0031] all involve flat, essentially rectangular panels, seeming to be 15 to 20″ in length, that extend rigidly straight back from the aft face of the trailer. This makes the devices susceptible to damage from accidental side or rear impact.\nU.S. Pat. App. No. 2002/0030384, Basford, employs drag vortex generators in conjunction with boattail plates and claims that the effective length of the boattail plates can thereby be decreased. A length of 12″ is described for these plates that are not self-stowing. Even this reduced length does not alleviate the potential for accidental damage.\nGerman Pat. No. DE2218300 discloses design features for corner contours to reduce aerodynamic drag on vehicles. These are intended to be formed on the actual exterior of the vehicle body and do not apply to added devices for heavy trailers."}
{"text": "The transport of truck trailers by train has traditionally been limited to long haul transport. Trailer design, railroad car design and scheduling problems have made it difficult for trains to compete in short and medium haul corridors. To date, trailers have had to be reinforced to handle the stress of rail transport. In addition, trucking companies have been forced to accept fairly wide windows of time for delivery.\nCurrently, trucking customers arrive at the train terminal and wait in line for a clerk to manually input the details of the shipment at the gate, to include billing information, contents of the truck, owner, driver, destination, weight, commodity code, commodity description, package type, broker and other pertinent information needed to transport the trailer and provide billing information. This process is time consuming and tedious; tractor-trailer drivers may wait in line at the gate for over an hour as each truck ahead of his goes through the check-in process. In addition, the time required to check-in a trailer is lengthy.\nFurthermore, transport of trailers via rail has been unpredictable. When the trucking customer arrives with a trailer to transport, they do not know if their trailer will be transported on the next train. The customer's container may be left behind for subsequent trains. In addition, it can be difficult to determine the train the trailer will be on and when that train will arrive at its destination. As a result trailer transportation via rail is an unreliable, frustrating and complex system.\nTrucks are a very efficient transportation medium, especially in short and medium haul corridors. Over the last 10 to 15 years, there have been a variety of attempts by the railway industry to move truck trailers onto rails. As noted above, most of these attempts to create intermodal transportation systems, however, require expensive modifications to the trailer to make it sturdy enough to withstand forces applied by cranes and by transport on the rails. Also the cost of building and outfitting railway terminals for truck transport is very expensive. A complex railway system for loading and loading trains is needed. In addition, trucking companies chafe at the long delays and complex processes of the railway companies.\nWhat is needed is a simple and reliable transport system which ensures that a train leaves on schedule and arrives at its destination within a predetermined, narrow window of time.\nIn addition, what is needed is a low cost high quality profitable intermodal system that reduces current bottlenecks and provides an efficient means of transporting trailers via rail.\nFinally, what is needed is a system and method for efficiently and economically transporting trailers via train across high density corridors."}
{"text": "The growth in personal computer (PC) usage and connectivity has forged new and creative ways for individuals and businesses to communicate and share information. Among these new methods include those which add video sequences to traditional methods of communication. Adding video sequences allow individuals and businesses to participate in activities such as video conferencing, video mail, interactive document editing, and object display.\nVideo sequences often contain large amounts of information. In order to transmit video sequences over bandlimited channels such as the Public Switched Telephone Network (PSTN), local area networks (LANs), the Internet, broadcast, and mobile networks, data rate control is required for encoding the video sequences. Data rate control may also be needed for encoding the video sequences for transmission over non-band-limited channels. Recording mediums such as hard drives, compact disk read only memory (CD-ROM) drives, and digital video disk (DVD) drives have storage constraints that need to be satisfied.\nA necessary parameter in the process of data rate control is the frame rate of the video sequence. In real-time environments, the exact frame rate is often non-deterministic. An example of this uncertainty in the frame rate is observed when fluctuations in system load on a computer system may temporarily increase or decrease the video capture frame rate. Many application program interfaces (APIs) such as Microsoft's Video for Windows.RTM. API fail to report the frame rate to installed video encoders during real-time capture and compression."}
{"text": "1. Technical Field\nThe present invention relates to prefetching in general, and more particularly, to a method for performing prefetching. Still more particularly, the present invention relates to a method for performing memory prefetching using an adaptive stream detection scheme and adaptive scheduling.\n2. Description of Related Art\nEarly prefetching techniques mainly focused on exploiting streaming workloads. While regular forms of spatial locality are relatively easy to predict, traditionally it has been difficult to exploit irregular patterns of spatial locality and it is even more difficult to exploit low amounts of spatial locality.\nRecently, a class of aggressive prefetching techniques has arisen from the notion of Spatial Locality Detections. Those techniques track accesses to regions of memory so that spatially correlated data can be prefetched together. The main advantage of those techniques is their ability to exploit irregular forms of spatial locality. However, the main disadvantage is their reliance on massive tables that tends to occupy large chip area and consumes a lot of power.\nConsequently, it would be desirable to provide an improved method for performing prefetching."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may or may not constitute prior art.\nThe reduction of noise, vibration and harshness (NVH) in passenger cars, light trucks and similar motor vehicles has always been a subject of concern and study. Recently, increased emphasis has been placed on the reduction of road noise, a major portion of which is generated by the rotating contact between the vehicle tires and the roadway. While attention has been directed to the tires and the tire/road interface, significant attention has also been directed to isolating the vehicle cabin from such exterior sounds with various materials and techniques.\nIn this regard, it has been found that significant noise pathways exist between the wheel wells, where road and tire noise is concentrated and the adjacent body structure such as the front foot wells."}
{"text": "The invention relates to a method for passive determination of target data associated with a target from measured and estimated bearing angles of the type mentioned in the precharacterizing clause of claim 1, to an apparatus for carrying out the method according to claim 10, and to a computer program which has suitable program code means for carrying out the method.\nIn order to determine the range, velocity and course of a target, for example of a surface vessel, a submarine or of a torpedo, as target data associated with the target from a carrier vehicle, for example a surface vessel or a submarine, without giving itself away, a sonar receiving installation is used to receive sound waves of the target noise, and to measure bearing angles to the target. These bearing angles are used to estimate a position of the target as a function of the carrier vehicle's own position, and to calculate an associated estimated bearing angle. This is dependent on the target moving uniformly, that is to say on a constant course without acceleration.\nDE 34 46 658 C2 discloses a method which iteratively reduces the difference between measured and estimated bearing angles over a plurality of processing cycles, until an error limit is undershot. When this error limit is undershot, the estimated position on which this is based is identified as the target position. The target data associated with this target position is therefore optimized on the basis of a predetermined optimization criterion, depending on the iteration method that is used. These include an optimized solution, which is updated with every new bearing angle measurement during each processing cycle. In uniform conditions, in particular without any target manoeuvre or change in the sound propagation, the optimized solution will evermore closely approach the actual solution as the observation time increases, because the solution which is in each case defined as the best solution, that is to say the optimized solution, converges with the actually correct solution at an earlier or later stage depending on the optimization method that is used and the current scenario.\nOne such method, which is also referred to as target motion analysis (TMA), has the disadvantage, however, that the calculated solution may be unrealistic because of unknown disturbances in the received signals. If the input data is received with major errors, the TMA method sometimes estimates unrealistic or illogical target data, such as a velocity of 80 knots or a departing course, even though the target is on an approaching course.\nWO 2009/156337 A1 discloses a method in which the reliability of the optimized solution is indicated in addition to this solution. For this purpose, a multiplicity of different target tracks and a quality measure for each of these assumed target tracks are calculated during each processing cycle. The reliability of the optimized solution can be deduced from the distribution of the quality measure.\nA further method is disclosed in DE 10 2007 019 444 B3, in which bearing angle measurements from a second transducer arrangement, which is arranged at a distance from a first transducer arrangement, are used for estimation of the target data when a parallax value, which is determined between the bearing angle measurements from the two transducer arrangements, is greater than a predeterminable parallax limit value and the range estimated for the target data estimate is less than a comparison range determined from the parallax limit value.\nEP 1 531 339 A2 describes a method in which limit values are predetermined for the target data, thus resulting in the iterative estimation process producing only physically and technically sensible target data. Completely unrealistic or illogical results, such as a range of 1000 kilometers, which can invariably occur in the iteration process when the received signals are subject to heavy noise, are excluded from the start.\nAn optimization method such as this, in which limit values are taken into account for the target data during the optimization process, is referred to as constraint TMA (CTMA). By way of example, the limit values are entered manually by an operator or else are determined from geographic characteristics. Even coarse presets make it possible to exclude estimates of target positions which could not possibly be detected by the sonar receiving installation.\nHowever, the known CTMA method has the disadvantage that its optimized solution cannot converge with the actually correct solution if the limit values are chosen incorrectly. Since the CTMA solution corresponds to the global minimum of an optimization criterion within the defined limit values, this does not produce the actually correct solution when the actual solution is outside the solution space defined by the limit values. After a convergence phase, the simple TMA method would therefore produce more accurate results than the CTMA method if the limit values had been chosen incorrectly.\nThe invention is therefore based on the problem of identifying and correcting incorrectly chosen limit values when using an optimization method for determination of target data."}
{"text": "1. Technical Field\nThe present invention is related to an apparatus and method for moving a first object relative to a second object and holding them in a defined position. In particular, the present invention relates to a scanning system for micromechanical devices.\n2. Description of the Prior Art\nWhen an object is to be moved or driven with a defined precision, then a device that controls the object's position and movement is employed. Such a device is a scanner that includes an actuator, whereby various kinds of driving principles are known for those actuators. Well known are electromagnetic, electrostatic, superconducting, piezoelectric, ultrasonic, pneumatic, air-based, thermal expansion, bimetal, and shape-memory alloy actuators. Micro-actuators and the micro machines driven by them are expected to have wide varieties of applications in information technologies, in medical, bioengineering, optics, and robotics fields. For example, electromagnetic actuators are ubiquitous. They can be found in everything from a large train to the smallest time piece. However, even the smallest magnetic actuators are usually made with wound coils and magnetic materials.\nIn the European patent application publication no. EP 0 998 019 A1, currently owned by the present applicant, a magnetic scanning or positioning system with at least two degrees of freedom is described. The magnetic scanning or positioning system comprises a supporting base equipped with magnets, a movable platform equipped with at least two electrical coils, and suspension elements providing an elastic connection between the movable platform and the supporting base. The electrical coils are positioned flat on the movable platform, thereby forming an essentially flat arrangement with the movable platform. The scanning or positioning with at least two degrees of freedom can be used in the field of scanning probe microscopy or in the field of data storage or imaging.\nAnother example is known as an electrostatic actuator. European patent application publication no. EP 0 865 151 A2 describes such an electrostatic actuator. The electrostatic actuator comprises a first member and a second member. The first member has a first opposed surface that includes an array of driven electrodes with pitch and the second member has a second opposed surface including an array of drive electrodes. A support positions the first member adjacent the second member with the first opposed surface spaced apart from the second opposed surface by a spacing. The ratio of the spacing and the pitch should be less than eight, and is preferably less than 2.25. The support permits the first member to move relative to the second member, or vice versa. A voltage source establishes a spatially substantially alternating voltage pattern on the array of driven electrodes. An electrode control establishes a substantially alternating voltage pattern on the array of drive electrode, and selectively imposes a local disruption on the substantially alternating voltage pattern on the array of drive electrodes to move the movable one of the first member and the second member relative to the other.\nEuropean patent EP 0 071 066 B1, granted to the present applicant, discloses an electric traveling support. The traveling support comprises a piezoelectric plate resting on three legs whose bottom surface is insulated from the bench on which the support is to travel, by a dielectric. The piezoelectric plate can be caused to contract by means of an actuating voltage applied via sliders to top and bottom electrodes on the plate. By applying a voltage to the legs, these may be clamped selectively by electrostatic forces effective across the dielectric. Appropriate control of the actuating and clamping voltages causes the support to either move in a linear or rotary fashion.\nAll the aforementioned actuators and scanners need energy not only for movement but also for holding and positioning. Low power precise positioning of parts and devices become more and more important as (nano)micro-mechanics gains in importance. In particular in nano-mechanics on the bases of local probes where tips have to be positioned with nano-meter precision reliable scanners are missing. For piezo-scanners domain wall creep and for electromagnetic scanners thermal creep are unsolved problems. Additionally, present scanners consume energy for holding and restoring the last position. Furthermore for conventional scanners the energy needed for moving from position to position increases with increasing spring deformation.\nIt is therefore an object of the present invention to overcome the disadvantages of the prior art.\nIt is another object of the present invention to provide a scanner with integrated tracking and nearly zero power consumption for holding a position with high precision."}
{"text": "1. Field\nThe disclosure relates generally to the data processing field and more specifically to migrating a source page containing a plurality of contiguous sub-pages that are subject to access by input/output devices, by means of a splitter, a selector, and an enhanced migration descriptor.\n2. Description of the Related Art\nReconfiguration of a computer system may be performed while the computer system is running. Doing so, however, frequently necessitates that data processing operations be disrupted. For example, if multiple operating systems are running on a computer, there may be a need to reallocate a block of memory being used by a first operating system so that it may be used by a second operating system. Traditionally, this situation has required that the first operating system stop using the block of physical memory before it is reallocated. As another example, when a problem is detected in a block of physical memory, it may be desirable to remove the memory from operation so it can be replaced. Again, the operating system that was using the block of memory must stop using it during the replacement process.\nA computer system may often be reconfigured without disrupting data processing operations by using virtual memory management mechanisms that are available in current operating systems. These mechanisms, however, are not satisfactory when the block of memory to be moved is subject to access by input/output (I/O) devices, commonly referred to as direct memory access (DMA). In particular, blocks of memory that are subject to access by I/O devices are conventionally “pinned,” that is, exempted from being moved, because access to them is maintained by essentially endless I/O programs that run as long as the operating system runs; and it would be disruptive to the system to stop such an I/O program.\nIt is possible to invalidate an I/O address translation entry for a block of memory (also referred to as a “page”) that is to be migrated. Then, if an I/O device were to attempt to access the page, the invalid translation entry would cause an I/O page fault to be reported and the I/O device driver would have to respond. That is, the driver would restart the I/O operation after the page again becomes available. In order for this procedure to work, however, the device drivers of all I/O devices that could possibly have access to the page would have to implement I/O page fault recovery code, which is not a practical solution.\nPrior mechanisms for migrating data subject to access by I/O devices have the problems discussed above. One solution involves a mechanism for migrating data in which DMA operations for I/O adapters that access the data being migrated are disabled for a short period of time. While this mechanism is effective in many applications, the industry standard PCI-E bus does not contain a mechanism for temporarily disabling DMA operations without adversely affecting other transfers to/from an I/O adapter.\nAnother solution, that allows migration of data without temporarily disabling DMA operations, is to use multiple migration descriptors. Methods to implement multiple descriptor functionality were originally disclosed in U.S. patent application Ser. No. 11/380,052. As will be explained further below, the method of U.S. patent application Ser. No. 11/380,052 requires the creation of a migration descriptor for each I/O mapped logical page that covers a real address range to be migrated.\nThus, when a block of data covers multiple real address ranges, a migration descriptor is required for each I/O mapped logical page covering a real address range to be migrated. In other words, a large page can have a number of sub-pages, each of the sub-pages having a different real address range. Presently, each of these source sub-pages requires its own migration descriptor in order to be migrated to a corresponding destination sub-page in a destination page. The allocation and management of multiple migration descriptors, where a migration descriptor is required for each I/O mapped logical page, introduces a significant overhead into the memory migration process. Therefore, a need exists for a way to migrate data without the memory load caused by multiple migration descriptors."}
{"text": "The present invention is directed to a catalytic process for selectively oxidizing carbonyl sulfide (COS) and carbon disulfide (CS2) to carbon monoxide (CO) and sulfur dioxide (SO2). The process of the invention is particularly useful in combination with the known equilibrium reaction for producing hydrogen from hydrogen sulfide. Such combination provides a cost effective way of producing hydrogen from what are often considered waste sulfur streams. The so-produced hydrogen can be used for a variety of purposes as is well known in the art.\nCarbonyl sulfide (COS) is found in many industrial process streams such as those associated with natural gas production, petroleum refineries and coal gasification plants. COS often is catalytically converted by hydrogenation to methane (CH4), hydrogen sulfide (H2S) and water; by hydrolysis to carbon dioxide (CO2) and H2S, or by oxidation to CO2 and sulfur dioxide (SO2). Hydrogenation is an expensive waste of valuable hydrogen (H2), while hydrolysis and total oxidation produces CO2, considered by many to be an environmental threat due to its alleged impact on global warming.\nCarbon disulfide (CS2) is another organic sulfur pollutant. CS2 is commonly encountered in paper manufacturing waste streams and in petroleum refineries. It can be treated much like COS to yield the same end products. Even though these methods provide the important removal of environmentally unsafe reduced sulfur compounds, they do not produce valuable chemicals from the original sulfur pollutants.\nH2S is commonly encountered in petroleum refineries where hydrodesulfurization is widely used to remove sulfur from gasoline. A significant amount of H2S (oftentimes as much as 10-30%) also is found in natural gas recovered from geological formations. The H2S is typically separated from the methane and the resulting de-sulfurized methane typically is distributed for industrial and personal uses, such as for home heating. The isolated H2S often is converted to elemental sulfur via the Claus Process. In the Claus Process, a first portion of the separated hydrogen sulfide is converted (oxidized) to sulfur dioxide (SO2) and water and the remaining portion of the hydrogen sulfide is reacted with the water-laden sulfur dioxide in the presence of a suitable catalyst to produce additional water and elemental sulfur. The so-produced sulfur represents a low value-added, commodity product, while essentially all of the potentially valuable hydrogen resident in the H2S is converted into water. The development of a cost effective alternative for using the hydrogen sulfide, and especially for recovering at least a portion of its hydrogen content, would significantly improve the economies of sour natural gas recovery and processing and would greatly benefit the operation of petroleum refinery operations.\nHerrington et al., U.S. Pat. No. 4,618,723 describes processes for using H2S as a reducing agent to reduce carbon oxides and generate hydrogen which is then used, as described, to make such organic compounds as alkanes, alcohols, alkenes and the like. The focus of such processes is to upgrade carbon dioxide and carbon monoxide to valuable organic compounds using H2S as a reducing agent. The patent presents the equilibrium reaction in which H2S is used to produce H2 from carbon monoxide, with by-product COS. The patent also shows the oxidation reaction involving COS to produce CO and SO2. This reaction, however, is conducted either in the absence of a catalyst or in the presence of quartz chips (SiO2) as a catalyst. Unfortunately, the yield (conversion) and selectivity of this step, as obtained under these conditions, has simply not been sufficiently high to make hydrogen production from H2S a cost effective alternative.\nThus, it would be advantageous if a process were available for selectively converting organic sulfur pollutants, such as COS and CS2, at high yield (conversion) and selectivity to CO and SO2. Such a process would make practicably feasible the conversion of H2S to the valuable end product, H2.\nU.S. Pat. No. 4,427,576 generally describes catalytically oxidizing H2S and/or organosulfur compounds to SO2using titanium dioxide (titania) with zirconia or silica (SiO2), an alkaline earth metal sulfate and one of the elements Cu, Ag, Zn, Cd, Y, lanthanides, V, Cr, Mo, W, Mn, Fe, Co, Rh, Ir, Ni, Pd, Pt, Sn and Bi.\nThe present invention is directed to a process for selectively oxidizing carbonyl sulfide (COS) and carbon disulfide (CS2) to carbon monoxide (CO) and sulfur dioxide (SO2):\nCOS+O2xe2x86x92CO+SO2\n2CS2+5O2xe2x86x922CO+4SO2\nThe CO can be used to produce valuable hydrogen, such as via the water gas shift reaction;\nCO+H2O⇄CO2+H2\nor more preferably from H2S via the known equilibrium reaction:\nH2S+CO⇄COS+H2\nIn this later reaction, the COS, in turn, is used to generate additional CO for the above reaction with H2S. The by-product SO2 also can be converted, if desired, to elemental sulfur, or to sulfuric acid.\nIn particular, the present invention is based on the discovery that by contacting COS, and/or CS2, in the presence of excess oxygen, with a supported metal oxide catalyst, one can selectively oxidize these organic sulfur compounds to CO and SO2. In particular, applicant has discovered that a supported metal oxide catalyst using the oxide of one or more of a variety of different catalytic metal oxides, i.e., comprising a metal oxide of a metal selected from the group consisting of vanadium (V), niobium (Nb), molybdenum (Mo), chromium (Cr), rhenium (Re), titanium (Ti), tungsten (W), manganese (Mn), tantalum (Ta) and mixtures thereof, supported on a metal oxide support selected from the group consisting of titania, antimony oxides, tantala, tin oxide, lanthana, indium oxides, iron oxides, nickel oxide, cobalt oxide, gallium oxides, manganese oxides, chromia, tungsten oxide, hafnia, zirconia, ceria, niobia, silica and alumina, such as vanadia on a titania support, is able to promote the selective oxidation of COS and CS2 to CO and SO2. As a general rule, the metal oxide support and the supported metal oxide surface layer should not be the same.\nOn the basis of cost and availability the preferred metal oxide supports are selected from the group consisting of titania (TiO2), zirconia (ZrO2), ceria (CeO2), niobia (Nb2O5), silica (SiO2) and alumina (Al2O3). Vanadia is often preferred as the metal oxide surface coating. Especially preferred catalysts are a 5% V2O5 surface coating on a Nb2O5 support (5% V2O5 on Nb2O5 catalyst), which for COS oxidation at 290xc2x0 C. provides a selectivity to CO of 98% and for CS2 oxidation provides an 80% selectivity to CO and a coating of V2O5 on a SiO2 support, which for CS2 oxidation provides a selectivity towards CO higher than 90% at 270xc2x0 C.\nBy coupling these oxidation reactions with the known equilibrium reaction involving the conversion of H2S and CO to H2 and COS (and often some CS2 and CH3SH as byproducts), one is able to convert what are otherwise environmentally undesirable sulfur pollutants to a very valuable material, H2 and by-product SO2. The SO2 is not laden with as much moisture as when produced directly from H2S oxidation and thus can be sent to a Claus plant for more efficient conversion to elemental sulfur, or alternately it can be further oxidized to sulfuric acid (H2SO4)."}
{"text": "Crawling, flying, and pestilent insects are troublesome to catch and dispose of, particularly for individuals who are afraid of, or repulsed by, such creatures. Existing devices such as fly swatters are currently used to crush such insects. However, the messy carcass of the dead insect is left remaining on the walls and/or on the fly swatter. Prior art fly swatters are known to have used frontal adhesives, but the frontal planes of the fly swatter do not provide a readily available means of rapidly, effectively, and hygienically entombing and completely sealing the flying and crawling insects. Moreover, such devices also fail to adequately seal the exposed adhesive surfaces on the device after insect contact, adhesion, and sealing. This is necessary to prevent future coincidental adhesive contact with surrounding surfaces or objects, including the user of the device.\nAnother disadvantage of prior art devices is that insect contact is established on an essentially random basis, versus the will of an operator on an immediate need basis. Additionally, use of chemical insecticides in human habitats may be undesirable and may be harmful, either short or long term. This is particularly true in nurseries, and/or in the presence of the very ill or elderly. The present invention eliminates the need for chemical insecticides in human habitats where their use could be potentially harmful.\nAdhesive technology exists today wherein the adhesive qualities can be modified to specifically address the user's requirements or desires. Additionally, pressure-sensitive adhesive papers or other materials are manufactured today at low unit-costs. Such papers might include the incorporation of complex or simple die-cuts and folds. Hence, what is needed is a device which utilizes such adhesive technology providing an operator means for selectively capturing and quickly disposing of an insect.\nStill another disadvantage of the prior art which uses adhesive is that stationary insect traps such as fly paper can accumulate large numbers of insects. Such accumulated pests might carry diseases or pathogens which could present a health-hazard potential. These accumulated insects, either dead or alive, are exposed to human surroundings for possible air borne or physical contact by humans and pets.\nA solution is needed which affords an operator a willful, portable means for quickly capturing and disposing of an insect in a hygienic and clean manner, even when an insect is at arms reach. This device can thereby incorporate an extension handle attached to a holder for an adhesive paper mounting device. Alternatively, a fixed location unit could also be provided which uses the same hygienic capturing and disposal system. The resulting portable or fixed device should allow for totally containing, sealing, and disposing of an insect in a fast, convenient and hygienic manner."}
{"text": "The following description relates to systems and techniques for reducing the polarization mode dependency of an optical device. The techniques may be used on in conjunction with optical devices such as those used in optical communications systems.\nDifferent optical devices, including optical amplifiers, optical filters and combined optical, amplifying filters may be used in optical communication systems. These devices may be integrated in, for example, components such as switches, routers, multiplexers and other components for optical-signal-processing.\nOptical signals that are employed in communication systems may be polarized and may include more than one polarization mode. Two polarization modes are known as transverse-electrical (TE) and transverse-magnetic (TM). The optical signals may propagate through the communication system on optical fibers and may have an arbitrary polarization state. The TM component may be thought of as propagating perpendicular to an axis of the optical waveguide and the TE polarization mode may be thought of as propagating parallel to the axis of the optical waveguide.\nThe performance of some optical components, including optical amplifiers and optical waveguide coupled components used in an optical communication system, may depend upon the polarization state of an incident optical signal received by the component. The optical devices may be polarization-dependent, meaning that the device has a different influence on the different modes of an incident signal.\nIncident signals with different polarization states may be affected by an optical device in different ways. In an optical amplifier, for example, the modal gain may be polarization mode dependent. The gain of the transverse-electrical and transverse-mechanical waves may be different. This variation in the response of an optical device to the different polarization states is known as polarization mode dispersion (PMD). Polarization mode dispersion is an inherent property of optical media. It can be caused by the difference in the propagation velocities of light in the orthogonal principal polarization states of the transmission medium. The net effect is that if an optical pulse contains both polarization components, then the different polarization components will travel at different speeds and arrive at different times, smearing the resultant optical signal. One result is that the gain may differ for TE-polarized and TM-polarized waves. The difference in gain between the differently polarized waves may result in an amplification of the TM wave that is different from the TE wave. Thus, the output optical signal from the amplifier may include TM and TE polarized waves that are in a different proportion than the input optical signal. The output signal of a polarization mode dependent device may have a different polarization state than the incident received signal."}
{"text": "Quantum optics now finds applications in quantum communication, quantum computing, and quantum sensing. In these quantum optics applications, photon sources are used that are designed to emit a photon stream having a controlled number of photons. In order to provide the necessary control over the output characteristics of the photon stream, these sources rely on materials and structures that are difficult to implement in chip-scale integrated circuit devices."}
{"text": "1. Technical Field\nThis invention generally relates to the field of projectiles, and more specifically relates to projectiles with deployable members.\n2. Background Art\nMany advances have been made in the art of projectiles, such as bullets fired from guns. Several known bullets are made of lead or other soft material that expands (known as “mushrooming”) when the bullet hits. The expansion of a lead bullet inside a target causes a greater knock-down effect, and increases the damage done to bones and internal organs, but typically slows the bullet to the point that it does not exit the target. In hunting applications, it is desirable for the bullet to exit the animal so the animal bleeds from the exit would, allowing the hunter to track the animal from the trail of blood. One way to assure the bullet exits the animal is to use a harder material that does not expand upon impact. The drawback of this approach is the damage done to the animal is not as great as for a softer, expanding bullet, increasing the likelihood of survival for an animal shot with a hard bullet. In addition, because a hard bullet does not expand, the animal will not likely bleed a great deal because the exit would is small, the same diameter of the bullet.\nSome projectiles have been developed with members that deploy to increase the damage when the projectile hits its target. For example, U.S. Pat. No. 6,240,849 to Holler and U.S. Pat. No. 1,464,032 to Daynix disclose projectiles that have members that deploy in-flight. These members increase the damage to the target upon impact. U.S. Pat. No. 1,318,858 to Frick discloses a projectile that may expand in-flight, or that may expand upon impact with a target. The Frick projectile includes pivoting knife arms that extend to create more damage to the target. The configuration of the Frick projectile is quite complex, and would be very difficult to manufacture in a cost-effective manner. What is needed is a projectile that provides members that deploy upon impact with a target that may be manufactured and assembled in a cost-effective manner."}
{"text": "The present invention relates to a cleaner for lithography which is used for dissolving or removing a resist etc. and in particular to a cleaner for lithography which can be used usefully to clean a coater such as in a coater cup, to remove an unnecessary resist from a substrate when or after coating of a resist, to remove the resist from the substrate after the object using the resist was achieved, and to clean or rinse the substrate after removal of the resist.\nFor production of integrated circuit elements, color filters, liquid crystal display elements, etc., lithography using resists are conventionally used. For example, in production of integrated circuit elements etc. by use of lithography, an anti-reflective coating is formed as necessary on a substrate, then a positive- or negative-working resist is applied thereto, the solvent is removed by baking, an anti-reflective coating is formed as necessary on the resist layer which is then exposed to ultraviolet rays, deep ultraviolet rays, electron beams, X-rays, etc., followed baking as necessary and development to form a resist pattern. Thereafter, baking is conducted as necessary, and the substrate is subjected to etching etc., and the resist is generally removed.\nFor coating of the resist, various methods such as spin coating method, roller coating method, reverse roller coating method, cast coating method, doctor coating method, dip coating method, etc. are used. For example, in production of integrated circuit elements, spin coating method is mainly used as a method of coating a resist. In spin coating method, a resist solution is dropped on a substrate, and due to the rotation of the substrate, this dropped resist solution is cast in the direction toward the outer periphery of the substrate so that an excess of the resist solution is removed by scattering from the outer periphery of the substrate, whereby a resist layer having a desired film thickness is formed. A part of the excess resist remains in a coater cup and solidified with time as the solvent is evaporated. The solidified material becomes fine powder which is then scattered to adhere to the substrate, thus causing resist pattern defects. To prevent this phenomenon, it is necessary to clean the coater cup every a few to dozens of substrates. Further, in the case where a resist layer having a desired thickness is formed by spin coating method on the substrate, there is the disadvantage that a part of the resist solution moves further to the back of the substrate, or the resist solution remains thickly (i.e. beads are formed) on the outer periphery of the substrate, and therefore it is necessary to remove the unnecessary resist or the beads from the side, the outer periphery or the back of the substrate. Further, even in other coating methods than spin coating method, adhesion of a resist to an improper portion occurs as is the case with spin coating. After the substrate is subjected to etching etc., the resist is generally removed and even in this removal step, the resist is dissolved and removed usually by means of an organic solvent.\nThe surface of the substrate from which the resist was removed by dissolution is made contaminant-free by generally cleaning using purified water etc., thus removing fine particles remaining on the surface of the substrate, followed by the subsequent step. If the solvent used for removal of the resist is a water-insoluble organic solvent or an amine type organic solvent, the substrate is often rinsed with a clean, water-soluble organic solvent rather than cleaning with purified water immediately after the step of removing the resist. The reason for this is that if a water-insoluble organic solvent is used as the liquid for removal of the resist, the resist dissolved in this solvent is prevented from adhering again to the substrate by preventing it from being precipitated in the solvent upon mixing with purified war, and that the water-insoluble organic solvent present on the surface of the substrate is once replaced by a water-soluble organic solvent to facilitate replacement by purified water. If an amine type organic solvent is used as the liquid for removal of the resist, this solvent remains in purified water to make it alkaline whereby the metallic substrate is prevented from corroding.\nAfter the coating step is finished, the coater should be cleaned for subsequent use or as a coater used again for other materials, or if an anti-reflective coating is resent between the substrate and the resist layer, the anti-reflective coating may be removed with a solvent as necessary after the resist pattern is formed.\nAs described above, a cleaner is used in various steps in lithography process, and as this cleaner, a wide variety of cleaners consisting of organic solvents are known (e.g. Japanese Patent Application Publication No. 4-49938).\nHowever, some conventional cleaners require much time or should be used in a large amount because of their poor ability to dissolve the resist and some are very toxic, so there is no cleaner satisfying both excellent solvent properties and safety to the human body, and therefore there is strong demand for a cleaner satisfying both solvent properties and safety to the human body. Although the production of integrated circuit elements has been described above, the above description applies not only to the production of integrated circuit elements but also to the production of color filters, liquid crystal elements, etc.\nThe object of the present invention is to provide a cleaner for lithography which is free of the aforementioned problems, that is, a cleaner which can be used usefully to clean a coater such as in a coater cup, to remove an unnecessary resist from a substrate when or after coating of a resist, to remove the resist from the substrate after the object using the resist was achieved, to clean and rinse the substrate after removal of the resist, which has a high solubilizing agility to permit sufficient cleaning to be effected rapidly in a small amount of the cleaner, and which is highly safe to he human body.\nAs a result of their eager research, the present inventor found that the above object can be achieved by using a homogeneous mixture of specific organic solvents and specific alcohols as a cleaner for lithography, to work out the present invention.\nThat is, the present invention relates to a cleaner for lithography which comprises a homogeneous solution comprising at least one selected from highly safe organic solvents of Group A below and at least one selected from alcohols having alkyl group with 1 to 4 carbon atoms.\nGroup A\npropylene glycol alkyl ether\npropylene glycol alkyl ether acetate\nethylene glycol alkyl ether\nethylene glycol alkyl ether acetate\nacetic acid alkyl ester\npropionic acid alkyl ester\nalkoxypropionic acid alkyl ester\nlactic acid alkyl ester\naliphatic ketone\nalkoxybutanol\nAnother present invention relates to a cleaner for lithography which comprises a homogeneous solution comprising propylene glycol alkyl ether and at least one alcohol selected from alcohols having alkyl group with 1 to 4 carbons.\nA further present invention relates to a cleaner for lithography which comprises a homogeneous solution comprising propylene glycol alkyl ether and at least one alcohol selected from alcohols having alkyl group with 1 to 4 carbons, wherein the propylene glycol alkyl ether is at least one selected from propylene glycol monomethyl ether, propylene glycol monoethyl ether and propylene glycol monopropyl ether, and the alcohol is at least one selected from ethanol, 1-propanol and 2-propanol.\nAn additional present invention relates to a cleaner for lithography which comprises a homogeneous solution comprising propylene glycol alkyl ether acetate and at least one alcohol selected from alcohols having alkyl group with 1 to 4 carbons.\nAnother further present invention relates to a cleaner for lithography which comprises a homogeneous solution comprising propylene glycol alkyl ether acetate and at least one alcohol selected from alcohols having alkyl group with 1 to 4 carbons, wherein the propylene glycol alkyl ether acetate is at least one selected from propylene glycol monomethyl ether acetate and propylene glycol monoethyl ether acetate, and the alcohol is at least one selected from ethanol, 1-propanol and 2-propanol.\nA still further present invention relates to a cleaner for lithography which comprises a homogeneous solution comprising at least one organic solvent selected from ethylene glycol ethyl ether acetate, n-butyl acetate, methoxypropionic acid methyl ester, ethoxypropionic acid ethyl ester, 2-heptanone, methoxybutanol and ethyl lactate, and at least one alcohol selected from ethanol, 1-propanol and 2-propanol.\nIn the present invention, preferable solvents in Group A are as follows:\n(1) propylene glycol alkyl ethers such as propylene glycol monomethyl ether, propylene glycol monoethyl ether, propylene glycol monopropyl ether, propylene glycol mono-n-butyl ether etc.;\n(2) propylene glycol alkyl ether acetates such as propylene glycol monomethyl ether acetate, propylene glycol monoethyl ether acetate etc.;\n(3) ethylene glycol alkyl ethers such as ethylene glycol monomethyl ether, ethylene glycol monoethyl ethylene glycol monopropyl ether etc.;\n(4) ethylene glycol alkyl ether acetates such as ethylene glycol monomethyl ether acetate, ethylene glycol monoethyl ether acetate etc.;\n(5) acetic acid alkyl esters such as propyl acetate, n-butyl acetate, n-amyl acetate etc.;\n(6) propionic acid alkyl esters such as methyl propionate, ethyl propionate, propyl propionate, butyl propionate etc.:\n(7) alkoxypropionic acid alkyl esters such as methoxypropionic acid methyl ester, ethoxypropinic acid ethyl ester, methoxypropionic acid ethyl ester, ethoxypropionic acid methyl ester etc.;\n(8) lactic acid alkyl esters such as methyl lactate, ethyl lactate etc.;\n(9) aliphatic ketones such as 2-butanone, 2-pentanone, 2-hexanone, 2-heptanone etc.;\n(10) alkoxybutanols such as methoxybutanol, ethoxybutanol etc.\nParticularly preferable solvents among those enumerated above are propylene glycol monomethyl ether, propylene glycol monoethyl ether, propylene glycol monopropyl ether, propylene glycol monomethyl ether acetate, propylene glycol monoethyl ether acetate, ethylene glycol monoethyl ether acetate, n-butyl acetate, methoxypropionic acid methyl ester, ethoxypropionic acid ethyl ester, ethyl lactate, 2-heptanone, and methoxybutanol.\nThe organic solvents in Group A can be used singly or in combination thereof. Preferable combinations of two or more organic solvents include a combination of a propylene glycol alkyl ether and a propylene glycol alkyl ether acetate, specifically a combination of propylene glycol monomethyl ether and propylene glycol monomethyl ether acetate.\nAlthough alcohols having alkyl group with 1 to 4 carbon atoms include methanol, ethanol, 1-propanol, 2-propanol, 1-butanol etc., ethanol, 1-propanol and 2-propanol are preferable in consideration of safety. These alcohols can be used singly or in combination thereof. Combinations of 2 or more alcohols include, for example, a combination of ethanol and 1-propanol and a combination of ethanol and 2-propanol.\nAlthough preferable examples of organic solvents in Group A or alcohols are illustrated above, the organic solvents in Group A or alcohols which can be used in the present invention are naturally not limited to those specifically illustrated.\nThe mixing ratio of the organic solvents in Group A and alcohols varies depending on which organic solvent or alcohol is used or which combination is used. However, the mixing ratio of the organic solvent in Group A: alcohol ranges from 10:90 to 99:1, more preferably from 40:60 to 90:10, most preferably from 60:40 to 80:20 by weight (hereinafter, xe2x80x9cratioxe2x80x9d is expressed on a weight basis).\nThe cleaner for lithography according to the present invention can be used in lithography using any positive- or negative-working resists known in the art. Typical examples of the resists to which the cleaner for lithography of the present invention can be applied include positive-working resists, for example, resists consisting of, e.g., a quinone diazide photo-sensitizer and an alkali-soluble resin and chemically amplified resists, and negative-working resists, for example, resists containing polymeric compounds having photosensitive groups such as polyvinyl cinnamate, resists containing aromatic azide compounds, resists containing azide compounds consisting of cyclized rubber and bis-azide compounds, resists containing diazo resin, photo-polymerizable compositions containing addition-polymerizable unsaturated compounds, chemically modified negative resists etc.\nThe above-described resist material consisting of a quinone diazide photo-sensitizer and an alkali-soluble resin is a preferable material to which the cleaner of the present invention is applied, and examples of the quinone diazide photo-sensitizer and alkali-soluble resin used in this resist material include quinone diazide photo-sensitizers such as 1,2-benzoquinonediazido-4-sulfonic acid, 1,2-naphthoquinone diazido-4-sulfonic acid, 1,2-naphtoquinonediazido-5-sulfonic acid, esters or amides of these sulfonic acids, and alkali-soluble resins such as polyvinyl phenol, polyvinyl alcohol, copolymers of acrylic acid or methacrylic acid, and novolak resin made e.g. from one or more phenols such as phenol, o-cresol, m-cresol, p-cresol, xylenol, etc. and aldehydes such as formaldehyde, paraformaldehyde, etc.\nThe chemically amplified resist is also a preferable resist to which the cleaner of the present invention is applied. The chemically amplified resist generates an acid upon exposure to radiation, and by chemical change due to the catalytic action of this acid, the solubility of the irradiated portion in a developing solution is altered to form a pattern, and for example, mention can be made of resists consisting of an acid-generating compound forming an acid upon exposure to radiation and an acid sensitive group-containing resin which is decomposed in the presence of an acid to form an alkali-soluble group such as phenolic hydroxyl group or carboxyl group, as well as resists consisting of an alkali-soluble resin, a cross-linking agent and an acid generator.\nThe cleaner for lithography according to the present invention can be applied to some compositions whose solvent is an organic solvent, out of known compositions for an anti-reflective coating which are applied for preventing a substrate from reflecting light upon exposure to light or for preventing interference of light from taking place in a resist layer.\nCoating of the cleaner of the present invention and the method of forming a resist pattern are specifically as follows: First, a resist solution is applied by a known coating method such as spin coating method to a silicon substrate, a glass substrate, etc. pretreated as necessary. An anti-reflective coating is applied and formed before or after coating of the resist as necessary. If this coating is formed by e.g. spin coating method, beads of the resist or of the anti-reflective coating tend to be formed on the edge of the substrate, but the fluidity of beads can be promoted by spraying the cleaner of the present invention onto the beads on the rotating edge so that a resist layer or an anti-reflective coating with substantially uniform thickness can be formed the substrate. In addition, the resist or the anti-reflective coating having moved to the outer periphery or the back of the substrate can be removed by spraying the cleaner to that portion.\nThe resist coated on the substrate is pre-baked on e.g. a hot plate to remove the solvent, thereby forming a resist layer usually having a thickness of about 1 to 2.5 xcexcm. Pre-baking is carried out at a temperature of usually 20 to 200xc2x0 C. and preferably 50 to 150xc2x0 C., although the temperature should be varied depending on the type of the solvent or resist used. Thereafter, the resist layer is exposed through a mask as necessary to radiation from known apparatuses for irradiation, such as high-pressure mercury lamp, metal halide lamp, ultra-high-pressure mercury lamp, KrF excimer laser, ArF excimer laser, soft X-ray irradiation apparatus, electron beam drawing apparatus etc., and after this radiation exposure, after-baking is conducted as necessary to improve development properties, resolution, pattern shape etc., followed by development to form a resist pattern. The development of the resist is conducted usually by using an alkaline developing solution. The alkaline developing solution used is a water or aqueous solution of sodium hydroxide, tetramethylammonium hydroxide (TMAH), etc.\nAfter the resist pattern is formed, the resist pattern is used for predetermined treatment such as etching, plating, ion diffusion, etc. for forming circuit elements, and the resist pattern is removed by use of the cleaner of the present invention. After this removal, the cleaner of the present invention can be used for rinsing as necessary. Further, in cases where the solvent for removal of the resist is a water-insoluble solvent or an amine type organic solvent different from that of the present invention, the cleaner of the present invention is used for rinsing as necessary, followed by rinsing with purified water. Further, the cleaner of the present invention can also be used effectively for cleaning of coaters."}
{"text": "Information related to a business, its customers, its competitors, and the like, is commonly referred to as “enterprise data” and can be used by the business for a variety of purposes. For example, enterprise resource planning (ERP) utilizes enterprise data to implement business functions such as manufacturing, supply chain management, financials, projects, human resources and customer relationship management (CRM). Enterprise data may also be used for product lifecycle management (PLM) to help manufacturers with product-related information. As a further example, enterprise data may be used for supply chain management (SCM) and supplier relationship management (SRM) functions.\nConventional software packages for providing access to enterprise data generally require significant storage and processing resources, and thus, are generally run on desktop computer systems, server computer systems, or the like. The computer system running the conventional software package may generate a query for specific enterprise data and send the query to a system able to access the enterprise data. Once data is received from the system, it may be processed by the computer system running the conventional software package and presented to the user of the computer system.\nAlthough enterprise data may be accessed using conventional software packages, the ability to access the enterprise data from remote locations or on the go is limited since the computer systems able to run the conventional software packages cannot be easily transported. Additionally, conventional software packages generally cannot be run on portable electronic devices such as smartphones or personal digital assistants (PDAs). Further, conventional portable electronic devices typically lack the storage and processing resources to adequately store and process the enterprise data. Thus, access to enterprise data is limited using conventional solutions.\nSome portable devices allow limited access to resident enterprise applications that run on the portable device. However, these portable devices limit communication between multiple enterprise applications. Unfortunately, enterprise solutions often require multiple enterprise applications to talk to one another to obtain information, etc. As such, portable devices that restrict or prevent communication between enterprise applications offer a poor platform choice for enterprise applications."}
{"text": "1. Field of the Invention\nThe present invention relates to tracking users in communication systems or networks supporting Broadcast-Multicast services (BCMCS).\n2. Related Art\nCurrent second generation and third generation wireless systems are primarily designed to support unicast voice and data services. The support of these unicast services to the end user has been achieved through various advances in wireless and networking technologies. The current emphasis in international standardization bodies such as 3GPP and 3GPP2 is on the design of protocols and procedures that allow the support of Broadcast-Multicast Services (BCMCS) over evolving networks. BCMCS is a bandwidth-conserving technology that reduces traffic by simultaneously delivering a single stream of information to a large number of recipients. Examples of these services include voice dispatch or Press-To-Talk (PTT) type services, broadcast/multicast streaming, etc.\nThere has been recent industry interest in the Public Safety Wireless Network (PSWN) context, where support of BCMCS may be especially important in terms of radio resource management. Radio resource management procedures may be needed to achieve low delay in establishment and delivery of BCMCS content, to provide scalability in terms of the number of multicast groups and/or number of users per multicast group, and to provide advanced service capabilities (e.g., ability to monitor/extract content from multiple BCMCS streams). Further, radio resource management procedures may be needed to achieve high spectral efficiency, to provide security (authentication, encryption and identity protection), and to improve mobile station battery life. Thus, radio resource management procedures that improve spectral efficiency and/or provide service flexibility are of interest to network operators.\nIn the unicast context, tracking of mobile stations within a geographic area is used primarily for the purposes of efficiently paging a mobile station. Mobile station tracking may perhaps become perhaps more important when assigning radio resources in the multicast context. In a cellular network, shared multicast radio channels should ideally be activated in a zone of cells or sectors for reliable reception. Depending on the level, or granularity, of tracking (i.e., accuracy of information at the network regarding the whereabouts of a BCMCS/PTT subscriber), radio resources may not be assigned in the proper sectors. As a result, there could be a significant loss in spectral efficiency for the network operator, and/or loss in quality of service perceived by a BCMCS/PTT subscriber.\nMobile stations that have an active uplink (e.g., those that send traffic or control signaling at regular intervals) can be tracked at the granularity of a sector. However, this is not expected to be the predominant operating mode for BCMCS/PTT where mobile stations are mostly idle or actively monitoring BCMCS content on the downlink. In such cases, it may be impractical for each user to maintain an active uplink, due to potentially large numbers of BCMCS subscribers. For such situations where maintenance of an active uplink is not feasible, tracking is typically accomplished through timer-based or event driven registration procedures. In addition to assisting procedures such as tracking, mobile station registration procedures also provide, to the network, information regarding subscription to specific services, radio priorities, etc. Timer-based procedures require mobiles to register with the network at given intervals (e.g. periodic intervals). Event driven triggers require mobiles to send a registration message to the network when an important event occurs (e.g. mobile decodes a change in location area).\nDepending on the options employed for registration, either fine tracking of mobiles down to the granularity of a sector, or coarse tracking of mobiles over a large number of sectors, is possible. If coarse tracking is employed, physical broadcast/multicast (BCMC) radio channels may need to be established in a large number of sectors that are not needed for reliable BCMCS content delivery, resulting in a waste of radio resources. Since it may not be possible to derive the benefits of power control in a network supporting BCMCS the power fractions needed to support a BCMC radio channel can be quite high. This limits the number of BCMC radio channels and/or unicast channels that can be simultaneously supported within a sector. Coarse tracking may therefore be unacceptable and may waste potentially scarce radio resources, as the network may have to transmit a BCMC stream within a large area, including sectors where mobile stations may not be present. On the other hand, tracking at the level of a sector (fine tracking) can result in excessive uplink overhead arising from frequent and sometimes unnecessary registrations by the mobile stations, thus overloading the uplink.\nMechanisms currently proposed for BCMCS registration include timer-based registration; registration on occurrence of certain events (e.g., power on, starting or stopping to monitor a multicast group); Registration Required Indication enabled in broadcast system parameter overhead on a new sector; and use of paging zone updates. These registration mechanisms are not well suited to tracking BCMCS subscribers that have any significant mobility (movement) across the network.\nWith timer based registrations, short timer values (the timer value may depend on the user mobility and desired accuracy) are needed for accurate tracking, otherwise the time elapsed from the previous registration may be quite long, potentially resulting in considerable ambiguity regarding mobile station location at the start of a BCMCS/PTT call. Long timers require radio resources to be assigned over a larger area to account for tracking uncertainty, thus resulting in spectral inefficiency. Although short timer values allow finer tracking, they result in large uplink overhead and reduce mobile battery life. They also result in repetitive unnecessary registrations from static (stationary) users.\nEvents such as “power on” and “start/stop monitoring” tend to be quite infrequent and dependent on end user behavior. As a result, these registration ‘triggers’ are not likely to be sufficient for the purpose of tracking. Relying on alternative fine tracking mechanisms such as the use of a “Registration Required” indicator bit within the broadcast system parameter overhead (i.e., group polling) of each sector may also give rise to a large number of spurious registrations from relatively static mobiles. Since the total number of BCMCS/PTT groups may be quite large, introducing such group polling mechanisms for all possible groups is inefficient.\nPaging zone updates have been proposed to reduce downlink paging load for packet data services in cdma2000®, Revision D. Each paging zone is assigned a paging zone identifier which is indicated in the broadcast overhead by all sectors within the paging zone. Mobile stations are associated with the paging zone corresponding to the sector having the strongest received pilot strength. The mobile stations report back to the network every time the paging zone changes. While this paging zone technique is reasonable in the unicast context, it does not allow efficient operation in the case of BCMCS. In particular, since mobile stations do not report updates regarding soft handoff candidates (i.e., sectors from which the received signal strength is sufficiently high) for a BCMC radio channel unless a sector from a different paging zone becomes the strongest in terms of pilot strength, this approach may not be suited to supporting soft handoff at paging zone boundaries. Moreover, hysteresis timers employed in paging zone updates for registration ensure that users will keep sending paging zone registration update messages upon hysteresis timer expiration, even if the paging zone corresponding to the strongest sector is unchanged."}
{"text": "In the standardization of 3GPP (Third Generation Partnership Project) there is on going work to standardize user authentication routines, especially for the so called Generic Authentication Architecture (GAA) involving a mutual authentication between a client and an application server. In the communication network several different applications will be available for the client and these applications will be supplied by third party suppliers, i.e. different from the Service Provider (SP). However, a client accessing several different applications would like to make only one single authentication, a so called Single Sign On (SSO) service, providing the possibility for the user to only authenticate once during a session, rather than to authenticate to each new application server it wants to use. This will make the authentication procedure much easier for the client. Also this kind of authentication service can be offered to third party application service providers as a service from the Service Provider handling the network. The GAA is aiming for solving this problem and making such a service available in the 3GPP network. The generic authentication provides an authentication of the users on an application level based on the proven security mechanism in the Public Land Mobile Network (PLMN).\nThe GAA is specified through the 3GPP group and drafts of the specification may be obtained through their web site, for instance the documents TS 32.220 and TR 33.919 maybe mentioned as good starting points regarding GAA. The GAA system may be explained as follows: A number of applications share a need for mutual authentication between a client/user (called UE, User Equipment, in the standard) and an Application Server (AS) in order to allow for further communication. This is necessary when the user wants to access servers demanding authentication, e.g. content servers charging for their services, certificate demanding web sites (e.g. banks), and similar application servers. Since many applications share the same need, it has been considered to specify a Generic Authentication Architecture (GAA), providing the architecture for allowing application servers access to the infrastructure authentication systems. Thus, if the application server trusts the service provider, this architecture may simplify the authentication schemes both for the user and for the application servers. The user needs only to authenticate once during a session, rather than authenticating towards every application server accessed.\nIn a GAA based session the user authenticates with the network infrastructure by providing an ID to a Bootstrap Function (BSF), this ID may for instance consist of the UE IMSI number (International Mobile Subscriber Identifier), which is a unique number coupled identifying a user. The IMSI is relayed to the Home Subscriber System (HSS, also called Home location register system) and the HSS provides an Authentication Vector (AV) to the BSF. The BSF authenticates the UE based on the USIM (Universal Subscriber Identity Module) and the UMTS-AKA methods, and sends a TID (Transaction Identifier) to the UE.\nAlso at the same time, work is on going to standardize so called Flow Based Charging (FBC). FBC has the aim to make it possible to charge users for service usage on a finer granularity than what is possible today. For instance it is of interest to identify the type of IP session that a user is running, the type of applications involved and so on. For instance one would like to be able to differentiate the charging costs for different types of services, e.g. streaming video may be charged more than exchanging plain text messages like simple email messages. There are many different services that may be used including both user to user and user to network services. Service data flows from these services may be identified and charged in many different ways. The FBC method is used to set up charging filters that is used by the CRF (Charging Rules Function) for different applications.\nThe filters provided in the FBC method can be quite complex and may involve source and destination address, source and destination port number, and transaction protocol, enabling a fine granularity of charging. Charging models requiring even more complex data may use special filters that look further into data packets and may be defined by the TPF (Traffic Plane Function) and invoked by the CRF.\nHowever, the above-mentioned two standardization works are not today aiming for an integration of the services they provide. This will in the future be crucial in order to be able to provide differentiated charging rules depending on the user and application connected to and at the same time ensuring the authenticity of the user towards the specific application server involved in the transaction. It should not be possible to, by only providing a fake IP number, getting access to services intended for other users or obtaining services at a wrong charging rate.\nThe work in 3GPP on GAA and FBC has been so far performed in parallel, and there is no concern taken to reuse functionality between the two functions. At some stage in standardization, interoperability between GAA and FBC needs to be built in. An integrated architecture will be necessary.\nSpecifically, the problem of supporting user specific charging for GAA authenticated users is not solved in FBC. The problem may be exemplified as follows:                1. A service provider has an application server providing a service, e.g. downloading a music file. This service is accessible for mobile users via the GPRS access network on the Gi interface. GAA ensures that the user is authenticated.        2. When the user downloads a music clip, he is charged by volume because the GGSN performs byte counting on the packets downloaded, and reports CDRs to the billing system. FBC can be used to enable that a specific rating is applied to the downloaded music file. FBC provides filtering in the GGSN that counts the specified service individually.        3. Now, the service provider might at some point provide the service for a special price for some users. In that case, it is needed to provide a specific filter for individual users. How that can be done is not specified today.        \nIn another standardization implementation, work is on going regarding policy decision functions such as the PDF (Policy Decision Function) for policy control of IP bearer resources, such as Quality of Service (QoS) for a specific user. Policy Decision Function relates a level of Quality of Service to a specific user and instant in time, enabling for instance a better Quality of Service for a certain type of application such as streaming video applications, or for a customer prepared to pay more for a higher quality of connectivity. It would be efficient to incorporate the PDF function with the GAA function in order to implement a more efficient network system."}
{"text": "1. Field of the Invention\nThis invention relates in general to the field of microelectronics, and more particularly to a microprocessor that provides a secure execution mode of operation that allows code to be executed in a highly secure environment within the microprocessor itself.\n2. Description of the Related Art\nThe use of desktop, laptop, and handheld computing and communication devices as platforms for digital communication of sensitive or proprietary data and digital rights controlled content continues to drive security innovations in the computing industry. For example, there are numerous applications which are disturbed at no cost over the Internet for the purposes of downloading and managing digital audio and video files. Through these applications, a user is allowed limited rights to songs, television shows, and movies. And a great deal of attention is paid to protecting those rights through the use of security features built into the applications which often rely on security mechanisms provided by their host platforms.\nIn addition to the protection of digital content rights, another factor which continues to drive improvements in computer system security is use restrictions which may be applied to the host platforms themselves. It is a well known fact that the cell phone industry has provided for so-called “pay-as-you-go” use of particular communication devices. With such a plan, a user pays no monthly fee, but instead pre-pays for a certain number of cellular minutes. When the minutes are exhausted, the user is denied access to the cellular network for any calls other than emergency calls.\nAs early as 2006, MICROSOFT® Corporation along with partner corporations, began to provide “pay-as-you-go” personal computing which is primarily directed toward emerging computer markets. Under this scheme, a user pays for their computer as it is used, through the purchase of prepaid cards. In addition, U.S. Patent Application Publication 20060282899, which is assigned to MICROSOFT, describes a system and method for delivery of an modular operating system, which includes a core function module, or basic kernel, providing fundamental operating system support and one or more add-on modules that allow customization of the operating system as desired. In this application, add-on modules may provide support or extended capability to the computer including hardware, applications, peripherals, and support. A digital signature may be used to confirm the integrity of an add-on module prior to installation, and certification may be verified to determine if installation of the add-on module is authorized. By withholding certification, a service provider may manage illegal or undesired modifications to a provided computer. In addition, digital rights management may be used to enforce terms of use of the add-on module in keeping with licensing arrangements.\nNot surprisingly, a veritable host of techniques have been developed as well which provide for circumventing the security measures that are now in place to protect and control access to rights controlled digital media, communication devices, and computing platforms. In more recent times, “hacking” has become a bona fide field of study. In fact, the present inventors have noted the publication of numerous works which are dedicated to tampering with or altogether defeating the security measures that have been put into place to safeguard access to and/or use of protected assets. One such work is the book, Hacking the Xbox: An Introduction to Reverse Engineering, by Andrew Huang, San Francisco: No Starch Press, 2003. The book focuses in particular on teaching hacking techniques to overcome the security mechanisms of the XBOX® gaming platform produced by MICROSOFT, and additionally provides significant teaching on the subjects of computer security and reverse engineering in general, with a discussion of the vulnerabilities of so-called “secure” computing platforms.\nConsequently, platform architects and designers continue to pursue techniques and mechanisms that are more effective at protecting against unauthorized platform access, regardless of whether the access is benign (e.g., probing, snooping), malicious (e.g., destructive or rights-defeating hacks), or somewhere in between (e.g., tampering). Many of these mechanisms are directed at preventing an intruder from physically accessing a platform, such as housing the platform in a secure chassis (e.g., a locked metal enclosure) or encapsulating vulnerable circuits in epoxy. But it is well known that these types of techniques add both cost and complexity to a system. Alternative mechanisms utilize security features that are provided for in a particular computer architecture itself.\nConsider the two primary security features provided for by the well-known x86 architecture: paged virtual memory and privileged execution. Under paged virtual memory, the underlying operating system defines a separate virtual address space along with access rights (e.g., execute only, read only) for each application that is being executed, thus precluding another surreptitious application from executing within the defined space or modifying its data. But the data associated with virtual address translation (i.e., page tables) can be easily snooped and changed since the data is resident in system memory and is presented external to the host microprocessor on its system bus.\nWith privileged execution, the architecture provides several levels of execution privileges, CPL0 through CPL3. Accordingly, certain system resources and instructions may only be accessed by applications which are executing at higher privilege levels. It is common to find operating system components operating at the highest privilege level, CPL0, and user applications relegated to the lowest privilege level, CPL3. But as one skilled in the art will appreciate, these architectural features were developed primarily to preclude system crashes due to software bugs, and are not very effective at preventing intentional or directed hacks.\nAs a result, various methods and apparatus have been developed which are more closely concentrated toward the prevention of intentional intrusion or takeover of a platform. In U.S. Pat. No. 5,615,263, Takahashi teaches a secure mode within a dual mode processor (i.e., microprocessor). In a general/external mode, the dual mode processor executes instructions provided from an external source. The instructions are supplied to the processor via input/output to the processor. Upon receiving a special software or hardware interrupt, the dual mode processor enters a secure/internal mode. The interrupt specifies a secure function stored in a read-only memory within the dual mode processor. Upon receiving such an interrupt, input/output to the dual mode processor is disabled. The identified secure function is executed by the processor. During execution of the secure function, any attempt to insert instructions not originating from the read-only memory are ignored. However, the processor may access data specifically identified by secure function being executed. Upon completion of performance of the secure function, an exit routine is executed to enable input/output to the processor and resume execution of instructions provided via input/output from the source external to the processor.\nTakahashi teaches that the secure mode is to be used for encryption and decryption and the dual mode processor relies upon normal instructions and data to be provided from an external control channel processor via a bus conforming to a standard bus architecture such as the Industry Standard Architecture (ISA). The dual mode processor powers up in non-secure mode and the secure mode is initiated via a software or a hardware interrupt. In secure mode, a limited number of functions (i.e., instructions) related to encryption and decryption can be executed. These functions are stored within a read-only memory (ROM) which is internal to the dual mode processor. As such, the present inventors have noted that Takashi's dual mode processor is inadequate in that it can only perform the limited number of functions which are provided for within it's internal ROM. Thus, an application comprising general purpose instructions cannot be executed in secure mode.\nIn U.S. Pat. No. 7,013,484, Ellison et al. teach a chipset for establishing a secure environment for an isolated execution mode by an isolated storage, which are accessible by at least one processor. The at least one processor has a plurality of threads and operates in normal execution mode or the isolated execution mode. The secure environment of Ellison et al. relies upon an external chipset (“isolated execution circuit”) which provides the mechanism for a processor to operate in isolated execution mode. The external chipset thus configures a secure memory area, it handles decoding and translation of isolated instructions, generation of isolated bus cycles, and generation of interrupts. While the external chipset indeed provides for proactive steps to isolate memory areas, instruction execution, and the like, it is noted that the external chipset is coupled to the at least one processor via a typical system bus, thus allowing for bus snooping and tampering with traffic on the bus itself during execution of any secure thread.\nIn U.S. Pat. No. 7,130,951, Christie et al. teach a method for controlling a secure execution mode-capable processor including a plurality of interrupts to interrupt the secure execution mode-capable processor when it is operating in a non-secure execution mode. The method includes disabling the plurality of interrupts from interrupting the secure execution mode-capable processor when it is operating in a secure execution mode. And while disabling interrupts is a desirable security feature in a secure execution environment, the processor according to Christie et al. relies upon instructions and data to be provided via an operating system over a system bus. Once the instructions are provided, then interrupts are disabled. Like the mechanism of Ellison et al., such an approach is clearly susceptible to bus snooping and tampering with any of the instructions which are passed to the processor over the bus.\nIn U.S. Pat. No. 6,983,374, Hashimoto et al. teach a tamper resistant microprocessor that saves context information for one program whose execution is to be interrupted, where the processor state is encrypted and stored to system memory. Hashimoto also teaches a technique for fetching encrypted instructions from system memory and apparatus for decrypting and executing the decrypted instructions. In addition, Hashimoto teaches using a symmetric key to provide the encrypted instructions in memory and then using an asymmetric key algorithm to encrypt the symmetric key, which is stored in memory. The symmetric key is known to the program creator and is encrypted using a public key that is read from the processor. The processor includes a unique private key that corresponds to the public key, which cannot be accessed by the user. Accordingly, upon execution of a branch instruction, program control is transferred to a “start encrypted execution” instruction which passes a pointer to the encrypted symmetric key. The processor fetches the encrypted symmetric key and decrypts it using its internal private key. Subsequently, the encrypted program instructions are fetched from system memory, decrypted using the decrypted symmetric key, and executed by the processor. If an interrupt or exception occurs, the state of the processor is symmetrically encrypted and saved to memory. Hashimoto teaches the use of common cache mechanisms, interrupt logic, and exception processing logic for both unencrypted and encrypted code.\nThe present inventors have noted that the microprocessor of Hashimoto is limiting in that the symmetric key corresponding to secure code is known by the code's creator, and could be compromised, thus exposing all systems having that code to security attacks. In addition, the present inventors have noticed that the processor of Hashimoto is disadvantageous in that decryption of secure code must be executed on the fly as instructions are fetched, which is extremely time consuming, thus causing throughput of the microprocessor to slow to a crawl. Furthermore, it is noted that the secure code of Hashimoto utilizes existing non-secure resources such as system memory, page tables, interrupt and exception mechanisms, all of which are subject to snooping.\nAccordingly, the present inventors have noted that it is clearly desirable to provide a microprocessor which is capable of executing an application or application thread comprising general purpose instructions (i.e., any of the instructions in the microprocessor's instruction set) within a secure execution environment.\nIt is additionally desirable that the secure execution environment be isolated from any of the known methods of snooping and tampering. Thus, it is required that instructions executed by a secure execution mode microprocessor be isolated from hardware within the microprocessor that provides access such as cache snoops, system bus traffic, interrupts, and debug and trace features.\nIt is furthermore desirable when applications are loaded for secure execution by the microprocessor, that a mechanism is provided to obfuscate the structure and content of the applications from any extant observation means and that a mechanism be provided to authenticate the source of the application and to confirm its veracity."}
{"text": "The present invention relates to an apparatus for helping to protect a vehicle occupant, and is particularly directed to an apparatus for helping to protect the lower extremities of a vehicle occupant in the event of a vehicle collision.\nIt is known to help protect a vehicle occupant from injury during a vehicle collision with an inflatable occupant restraint such as an air bag. Air bags are known to deploy into the passenger compartment of the vehicle from a variety of vehicle locations including the steering wheel, the instrument panel, the side doors, and the headliner, to help restrain movement of the vehicle occupant\"\"s head and torso. It is also known to position an inflatable restraint in the area of the knee bolster in the vehicle to help protect the knees of the vehicle occupant during a collision.\nA vehicle collision, particularly a severe frontal collision, can result in deformed portions of the vehicle body penetrating the footwell of the vehicle where the occupant\"\"s feet are located. Accordingly, an apparatus which moves the occupant\"\"s feet rearward in the footwell away from the penetrating vehicle body portions would help to protect the feet from injury during such a collision.\nThe present invention is an apparatus for helping to protect an occupant of a vehicle. The apparatus comprises an inflatable device having a stored deflated position in which the inflatable device is located adjacent an instrument panel in the vehicle, and an inflator for providing inflation fluid for inflating the inflatable device. The inflatable device, when inflated, engages at least one lower extremity of the occupant. An actuator moves the inflatable device rearwardly in the vehicle to move a foot associated with the at least one lower extremity of the occupant rearwardly in the vehicle."}
{"text": "The present invention relates to a coupling or latching assembly for removably coupling separate parts together, such as cover or panel or housing parts on a vehicle.\nCover panels on agricultural equipment have traditionally been attached with fastening hardware or latches which must be attached to the panel to facilitate mounting the panel on the vehicle. Such hardware is expensive, and labor is required to install such parts on the vehicle and subassemblies. When a cover is to be installed on a vehicle, such hardware must be carefully aligned so that the cover can be mounted correctly and retained securely during field operations. In some cases, the installation and function of the latch hardware may be compromised if the panels and mating parts do not have uniform wall thickness at the attaching points."}
{"text": "1. Field of the Invention\nThe present invention relates to conservation of electrical energy consumption by commercial, industrial, residential and all other energy consumers using retrofitted control devices which are attached at the incoming breakers to a facility and operate to increase efficiencies in a three stage process. Additionally, the present invention system includes lightning suppression and transient voltage surge suppression (TVSS) features, and failure detection circuitry.\n2. Information Disclosure Statement\nThe following references are examples of the prior art relating to control of electrical energy consumption:\nU.S. Pat. No. 4,163,218 relates to an electronic control system for controlling the operation of a plurality of electrical devices which are energized from AC power lines which includes a single, central unit connected to the power lines, which further includes a central transceiver means for transmitting an encoded oscillating signal of one frequency onto the power lines, a central encoding means for encoding the oscillating signal with an encoded signal in synchronization with the frequency of the AC power for selective control of electrical devices, and a central control means connected to the encoding means for selecting the electrical device to be controlled and its desired state. The invention further includes unitary switch units respectively interconnected between power lines and each electrical device being operative for both local and centralized control of the electrical device with the local control and the centralized control placing the electrical device in respective opposite states from each other, each switch unit including a switch transceiver means for receiving the encoded oscillating signal from the power lines, a switch decoding means coupled to the switch transceiver means for detecting the encoded signal, a switch control means connected to the switch decoding means for setting the selected electrical devices to the desired state, and a local control means for selectively locally operating the electrical device independently of the central unit and placing the electrical device in a state opposite from that which it was placed by the central unit.\nU.S. Pat. No. 4,845,580 describes a spike elimination circuit for A.C. and D.C. power sources which comprises two gas tube and/or two semiconductor voltage limiting devices before a Bandpass Filter. The Bandpass Filter consists of 2 capacitors to ground an inductor in series with the line. The spike eliminator can be portable, mobile, or hard wired for the protection of home controls and electronics, telecommunications, commercial and industrial controls and the computer field and others.\nU.S. Pat. No. 4,870,528 describes a surge suppressor which comprises a first series circuit having a first inductance and a first alternating voltage limiter, including at least a first capacitance and a bidirectionally conductive rectifying circuit for charging the first capacitance, coupled between first and second input terminals for limiting surge currents and voltage excursions coupled to first and second load output terminals. The first alternating voltage limiter further comprises a sensing circuit for sensing at least one of the charging current supplied to and the voltage developed across the first capacitance. An auxiliary energy storage circuit and a normally open switching device responsive to the sensing circuit are provided for coupling the auxiliary energy storage circuit across the first capacitance during high energy surge conditions.\nU.S. Pat. No. 5,105,327 describes a power conditioner for AC power lines which has a choke and capacitor coupled in series across the power lines. The choke comprises a coil terminating in a line, with the line looped back through the coil. The power lines are thereby balanced to provide greater operating efficiency. Capacitors and transient suppressors (e.g., varistors) are used for transient suppression and power factor correction.\nU.S. Pat. No. 5,420,741 relates to an arrangement for obtaining flux rate information in a magnetic circuit including passive means connected across a flux rate sensor for implementing control of said flux rate. The passive means being a tuned magnetic flux rate feedback sensing and control arrangement wherein impedance is tuned and the energy loss characteristic is adjustable. The selection of inductance and capacitance values provides tuning and the selection of resistance affects the energy loss characteristics.\nU.S. Pat. No. 5,432,710 is directed to an energy supply system for supplying, in system interconnection, power at a power receiving equipment from a power plant and power generated by a fuel cell to a power consuming installation, and supplying heat generated by the fuel cell to a heat consuming installation. This system includes an operation amount computing device for computing an amount of operation of the fuel cell to minimize an equation y=aXL+bXM+cXN, in response to an energy demand of the power consuming installation and heat consuming installation. A control device controls the fuel cell to satisfy the amount of the operation computed. The system supplies energy in optimal conditions with respect to the cost borne by an energy consumer, consumption of primary energy, and release of environmental pollutants. Energy is effectively used from the standpoint of the energy consumer and a national point of view.\nU.S. Pat. No. 5,436,513 relates to an information handling system which is described as having a power supply and having a switching circuit that switches a plurality of energy sources between series and parallel couplings. Associated with the switching circuit is a voltage level detecting circuit for monitoring the voltage level of the energy sources. A processor for controlling the information handling system responds to the voltage level detecting circuit and in the event of a low voltage condition, the processor activates the switching circuit to switch the energy sources and from a series to a parallel coupling. Alternatively, the processor responds to other inputs or conditions for actuating the switching circuit.\nU.S. Pat. No. 5,459,459 is directed to an algorithm for implementation in a meter register and a reading device. In the one embodiment, the invention enables selecting a display table to be read from the register, updating the billing read date and time in the register, reversing the order in which load profile data is transmitted from the register to the reader, specifying the number of load profile intervals to be read from the register, and specifying the number of intervals to skip when reading from the register.\nU.S. Pat. No. 5,462,225 relates to an apparatus and method for controlling energy supplied to a space conditioning load and for overriding a load control operation in response to measuring certain space temperatures within a closed environment. The load control apparatus includes a control device connected to an electrical distribution network and to a space/conditioning load and a temperature sensing device connected to the control device. The control device conducts a load shedding operation to control distribution of electrical energy to the space conditioning load in response to command signals supplied by a remote command center. The temperature sensing device operates to override the load shedding operation by outputting a control override signal to the control device in response to sensing certain space temperatures within the closed environment. If the temperature control device is connected to an air conditioning system, the temperature sensing device causes the control device to terminate the load shedding operation prior to expiration of a selected time period in response to measuring a space temperature that exceeds a maximum space temperature limit. In contrast, if the temperature control device is connected to a forced air heating system, the temperature sensing device causes the control device to terminate the load shedding operation when a measured space temperature drops below a minimum space temperature limit. The maximum space temperature limit is greater than the control temperature setpoint of a thermostat that controls the space conditioning operations, whereas the minimum space temperature limit is less than the control temperature setpoint.\nU.S. Pat. No. 5,483,672 relates to a communication system, where a communication unit may conserve source energy when it is inactive in the following manner. The control channel is partitioned into a predetermined number of windows and a system window which are transmitted on the control channel in a round robin manner. When the communication unit registers with the communication system, it is assigned to a window group. The communication unit then monitors only the system window to determine whether the window group that its been assigned to is also assigned to one of the predetermined number of windows. When the window that has been assigned to the window group is being transmitted to the control channel, the communication unit activates to monitor that window. Once the window is no longer being transmitted, the communication unit deactivates until the system window is being transmitted or the window assigned to the window group is being transmitted.\nU.S. Pat. No. 5,495,129 relates to an electronic device for multiplexing several loads to the terminals of a source of alternating electrical energy. The source of alternating electrical energy is coupled by electromagnetic flux to the loads by using primary excitation windings and connects to the terminals of the source of alternating electrical energy and secondary windings respectively corresponding to the number of loads. The secondary windings are at least partially coupled to the primary winding and are each connected to the terminals of a load. The coupling is inhibited by auxiliary winding which are each totally coupled with the secondary winding. The inhibition function is controlled in order to inhibit all the magnetic couplings except for one and this particular one changes as a function of the respective load to be coupled to the source of alternating electrical energy.\nU.S. Pat. No. 5,512,831 relates to a system for testing electrochemical energy conversion and storage devices includes means for sensing the current from the storage device and varying the load across the storage device in response to the current sensed. The system is equally adaptable to batteries and fuel cells. Means is also provided to sense system parameters from a plurality of locations within the system. Certain parameters are then stored in digital form for archive purposes and certain other parameters are used to develop control signals in a host processor.\nU.S. Pat. No. 5,517,188 is directed to a programmable identification apparatus, and associated method, includes a transceiver and a transponder. The transponder is powered by the energy of a transceiver transmit signal generated by the transceiver and includes a programmable memory element. A coded sequence which uniquely identifies the transponder is stored in the programmable memory element and, when the transponder is powered, the transponder generates a transponder signal which includes the coded sequence stored in the programmable memory element, once modulated by circuitry of the transponder. When the transceiver transmit signal generated by the transceiver circuitry is of certain signal characteristics, the coded sequence stored in the programmable element is erased, and a substitute coded sequence, which also forms a portion of the transceiver transmit signal, is stored in the programmable memory element. Storage of the coded sequence in the programmable memory element is, hence, effectuated merely by application of a transceiver transmit signal of certain characteristics to the transponder.\nU.S. Pat. No. 5,528,123 measures the total line current in a power cord which is used to energize both a power factor corrected system and non-power factor corrected AC loads. The power factor control loop of the power factor corrected system is then driven to correct the power factor of total line current in the power cord ideally to approach unity.\nU.S. Pat. No. 5,640,314 relates to a symmetrical ac power system which provides a balanced ac output, whose maximum voltage with respect to a reference ground potential is one-half the ac output voltage, and which is derived from a single phase ac source through the use of an isolation transformer having a center-tapped secondary winding. The center tap is connected to the output power load circuit as a ground reference potential with respect to the symmetrical ac output so as to constitute the reference ground potential for the power supply and load. Since symmetrical ac power is applied to the load by the system, reactive load currents, other power artifacts, EMI and RFI emissions and other interference and noise components ordinarily resulting from the application of conventional ac power to the load are reduced or eliminated by appearing as equal inversely phased signal elements which cancel one another. In order to maximize the performance of the symmetrical power system, the isolation transformer has a bifilar-wound secondary winding.\nU.S. Pat. No. 5,646,458 describes a UPS (uninterruptible power system) which includes an UPS power conditioning unit that provides conditioned AC power to a critical load. The UPS power conditioning unit includes a variable speed drive that operates in response to AC utility power or to a standby DC input by providing a motor drive signal. The UPS power conditioning unit further includes a motor-generator that operates in response to the motor drive output by providing the conditioned AC power to the critical load. In response to an outage in the utility AC power, standby DC power is provided by a standby DC power source that includes a variable speed drive and a flywheel motor-generator connected to the variable speed drive. Both the UPS power conditioning unit and the standby DC power source are initially operated in response to the utility AC power, the flywheel motor-generator storing kinetic energy in a rotating flywheel. When an outage occurs, the rotating flywheel continues to operate the flywheel motor-generator of the standby DC power source, causing the production of AC power which is rectified and provided as standby DC power to operate the variable speed drive of the UPS power conditioning unit until either the utility AC power outage is over or a standby emergency generator is brought on line.\nU.S. Pat. No. 5,880,677 relates to a system that monitors and controls electrical power consumption that will be retrofitted to a typical consumer electrical power arrangement (typical arrangement-electrical feed line from a provider, a meter, a circuit breaker and individual input wiring to a plurality of electrical devices, appliances and outlets). The system includes a control unit which receives information from an electromagnetic pickup device from which real time electrical consumption is determined over very short periods of time. The control unit has a main data processing and storage processor for retaining information and it may include a communication microprocessor for sending signals to corresponding modules. The electromagnetic pickup device uniquely measures the electromagnetic flux emanating at each output wire from each of the individual circuit breakers in a breaker box. The modules have filters which release electrical power to the individual electrical devices, appliances and outlets at a controlled, economic rate.\nU.S. Pat. No. 5,892,667 describes a symmetrical ac power system which provides a balanced ac output, whose maximum voltage with respect to a reference ground potential is one-half the ac output voltage, and which is derived from a single phase ac source through the use of an isolation transformer having a center-tapped secondary winding. The center tap is connected to the output power load circuit as a ground reference potential with respect to the symmetrical ac output so as to constitute the reference ground potential for the power supply and load. Since symmetrical ac power is applied to the load by the system, reactive load currents, other power artifacts, EMI and RFI emissions and other interference ad noise components ordinarily resulting from the application of conventional ac power to the load are reduced or eliminated by appearing as equal inversely phased signal elements which cancel one another. In order to maximize the performance of the symmetrical power system, the isolation transformer has a bifilar-wound secondary winding.\nU.S. Pat. No. 6,009,004 discloses a new single-phase passive harmonic filter for one or more nonlinear loads. The filter improves the total system performance by drastically reducing the line side current harmonics generated by nonlinear loads. The filter includes two inductive portions across one of which is connected a tuning capacitor. The parallel combination of one inductive portion with the tuning capacitor forms a series tuned filter configuration while the second inductive portion is used for harmonic attenuation. A shunt capacitor is employed for shunting higher order harmonic components. A single-phase passive voltage regulator provides the needed voltage bucking to prevent over voltage at the load terminals of the filter. The filter provides an alternate path for the harmonic currents generated by nonlinear loads. The over voltage caused by the increased capacitive reactance is controlled by either capacitor switching or by the use of the passive voltage regulator or a combination of the two. Capacitor switching is dependent upon load conditions.\nU.S. Pat. No. 6,014,017 describes a method and an apparatus for power factor correction for a non-ideal load, which is supplied from a mains power supply, by a compensation device which is electrically connected in parallel with the load and has a pulse converter with at least one capacitive store. A transfer function space vector is calculated as a function of a determined mains power supply voltage space vector, a mains power supply current space vector, a compensator current space vector and of an intermediate circuit voltage which is present on the capacitive store. As a result of which the pulse converter generates a compensator voltage space vector on the mains power supply side as a function of the intermediate circuit voltage. A compensator current space vector, that keeps the undesirable reactive current elements away from the mains power supply, is thus obtained via a coupling filter that is represented as a compensator inductance.\nU.S. Pat. No. 6,058,035 describes a method wherein after starting the input of a switching signal to a booster circuit whose boosting rate is changeable in accordance with the duty ratio of the inputted switching signal and calculating the output power of an invertor circuit, which is connected to the subsequent stage of the booster circuit, from the output current of the invertor circuit, the target voltage after boosting by the booster circuit is obtained based on the output power. If the actual output voltage of the booster circuit is lower than the target voltage, the duty ratio of the above switching signal is increased, and if higher, the duty ratio of the above switching signal is decreased.\nNotwithstanding the above prior art, there are no teachings or suggestions that would render the present invention anticipated or obvious.\nThe present invention is a power pod controller system, which includes connection for connecting the device in parallel with an incoming power supply to a facility, i.e. a home, factory, office, institution, etc.; there is an electricity pod controller device, which includes a first stage component, a second state component and a third stage component.\nThe first stage is adapted to recognize electromagnetic interference, and to respond thereto by suppressing live transient voltage surges, and thus includes at least one line transient voltage surge suppressor (TVSS). The TVSS may be one or more metal oxide varistors and is preferably a plurality of in-parallel metal oxide varistors.\nThe second stage component includes at least one variable inductor to regulate the total harmonics distortion (THD) and thus enhance power factor correction. This would include chokes (transformers) and storage. The third stage component includes hardware to provide power to maintain phase regulation to incoming power, thus, to maintain true phase relationship between voltage and current at times of increased power demands.\nThe third stage component includes at least one power storage and discharge element.\nIn addition, there is failure indicator circuitry connected to the electricity pod controller device and specifically to the surge suppressor(s) for indicating failures and for setting off appropriate responses, e.g. alarms, back ups, etc. at remote locations.\nIn some embodiments, the present invention is arranged so that the first stage component, second stage component and third stage component operate in a single phase. In other embodiments, there are at least two sets of components, each set having a first stage component, a second stage component and a third stage component, such that the device operates as a two phase device. In yet other embodiments, the present invention system with the electricity pod controller device has three sets of components, each set having a first stage component, a second stage component and a third stage component, wherein the device operates as a three phase device.\nMore specifically, preferred embodiments of the present invention power pod controller system include:\n(a) connecting means for connection to an incoming power supply of a facility, for connection in parallel, including a hot line and a neutral line, and at least one ground, with a number of electricity pod controller device components connected between the hot line and the neutral line. They are connected in the following order travelling in a direction away from the connections;\n(b) at least one front metal oxide varistor line transient voltage surge suppressor having a predetermined capability to suppress undesired power spikes (xe2x80x9cfrontxe2x80x9d is used to mean closer to the connection and xe2x80x9cbackxe2x80x9d is used to mean further away from the connection);\n(c) at least one capacitor of predetermined capacitance;\n(d) at least two chokes in the form of inductor/metal oxide varistor transformers;\n(e) at least a second capacitor of its own predetermined capacitance;\n(f) at least one back metal oxide varistor having a predetermined capability. In preferred embodiments, the front metal oxide varistor is a plurality of varistors in parallel. It has been discovered that a plurality of smaller varistors having the same total capability as one larger varistor responds more quickly than one large varistor. For similar reasons, it is preferred that the back metal oxide varistor be a plurality of varistors in parallel. However, single large varistors may be used without exceeding the scope of the present invention.\nIn addition, there is failure indication circuitry which is connected in series with the TVSS(s), which function as supervision and alarm devices. In some preferred embodiments, a plurality of lightning arresters are likewise connected to the TVSS(s).\nIn some embodiments, the components (b) through (f) above are arranged for operation as a single phase device, while in other embodiments, the components are duplicated therein to form two connected sets thereof and are arranged for operation as a two phase device. In yet other embodiments, at least a portion of component (e) is duplicated.\nIn those embodiments where three phase operation is desired, the electricity pod controller device components (b) through (d) are at least triplicated therein and (e) and (f) are at least duplicated therein to form three connected sets thereof and are thus arranged for operation as a three phase device.\nThe system of the present invention may further include at least one indicator lamp connection and lamp at the device itself, which is wired so as to illuminate when the device is functional and to not illuminate when the device is non-functional. This optional feature is different from the aforesaid failure indicators which are connected directly with the TVSS(s) and are usually used for remote signalling, such as a supervisor\"\"s workstation, or to a host server which will broadcast the alert signal to a plurality of locations to simultaneously warn supervisors, managers and even smart response systems."}
{"text": "This invention relates to the field of computer systems. More particularly, a system and methods are provided for resynchronizing a media stream to a client.\nMedia streaming systems are designed to stream media programs and events, which may be pre-recorded or live, to client devices (e.g., computers, media players). A client can play a media stream as it is received (i.e., before the stream is complete), thereby allowing fast, or even real-time, enjoyment of the media program or event.\nStreamed media is often composed of multiple tracks. For example, an audiovisual program generally includes at least one audio track and at least one video track. Each track includes media of the appropriate type (e.g., audio, video), and may also include metadata that is used to stream the media correctly. A track's metadata may include information for identifying a media segment or sample (or other unit of media) that should be played for a given time index within the program, for determining where that segment or sample is located in the file, etc. A media track's metadata may thus be used by a media streaming server to stream the media track in the correct sequence, with appropriate timing, and so on.\nWhen streaming a media track or program to multiple clients, existing systems usually store separate copies of each track's metadata for each client. More specifically, from a media file containing the media to be streamed, these systems repeatedly extract and store the metadata each time a new client requests the media. When serving a large number of clients, this could require a significant amount of the server's resources (e.g., memory, processor, disk usage), and may limit the number of clients the server can support.\nAlso, some existing systems typically allocate just one file descriptor, to be shared among all clients receiving a particular media stream, for accessing the media file containing the media being streamed. This can lead to a great deal of contention among the client streams as each one attempts to seek to (i.e., find) and extract a different media segment or sample.\nYet further, when streaming a multi-track media program to a client, existing media streaming systems attempt to maintain synchronization between the tracks so that the appropriate corresponding track media is streamed or played for each time unit of the program, but may be incapable of regaining synchronization if it is lost. For example, if the media for one track is received (e.g., from a storage device) at a slower rate than the media for another track, an existing media streaming system may simply continue streaming the media, even if the tracks get further and further out of synchronization. Other systems may simply halt a media stream if synchronization is lost, without attempting to correct the situation."}
{"text": "The invention relates generally to Web page filtering, and more particularly to a system and method by which end users may selectively filter out sections of Web pages, such as banner ads and/or other material.\nA vast amount of information is readily and publicly accessible to users on the Internet, and that collection is growing. Users may access the Internet to retrieve information from a great variety of Web pages, such as news from a news provider, search results from a search engine, product or service information, or other information. As users retrieve a desired Web page, often times other unnecessary information is automatically retrieved and displayed within the borders of that page. Such unwanted information may include advertisements, banners, animated links and other peripheral material. Although many Web site operators rely on the xe2x80x9ceyeballsxe2x80x9d they attract for advertising returns, the unchecked retrieval of advertisements may slow own the browser and hinder the user\"\"s Internet experience. Generally, advertisements are displayed on designated areas of the Web page, such as the top portion or sides of the page. Therefore, it is often easy to predict where advertisements will be positioned with respect to the rest of the Web page content. There may be instances where inappropriate advertisements or other material (e.g., adult themed promotions or unsecured ordering sites) may also be inadvertently retrieved and viewed by an unintended audience, such as young users.\nCurrent browsers offer a limited ability to block the download of all images in a Web page, after which the user must click on each image text tag to display the corresponding portion. It would be desirable to define sections of a Web page where download and display should be avoided, while other portions of the Web page are processed and displayed normally. One application would include blocking ad banners and other unwanted material, while leaving the rest of a Web page intact.\nCurrent browsers thus offer an all or nothing approach to filtering image downloads. It is difficult to define areas of a Web page which should be excluded from display, while all other Web page data is displayed normally. Other filtering tools exist which attempt to specifically filter objects which look like ad banners, based on their text description tag or image size and dimensions. Currently, users do not have the ability to define their own desired filter criteria which can block the download and display of any Web page content (image, text, or other content). Therefore, there is a need for an efficient and effective system and method for selectively filtering out unwanted material on Web pages, including so that the filters associated with the URL are applied each time the Web page is subsequently fetched.\nThe invention is directed to a system and method for selectively filtering out sections of a Web Page, such as banner ads or inappropriate material, from being downloaded according to a personalized, specific filter to omit certain text portions or areas of Web pages each time that site is revisited. The user may select one or more sections to be omitted by highlighting one or more items and capturing one or more associated text streams. Alternatively, the user may select the one or more sections to be omitted by drawing one or more boundary boxes around each item to be omitted.\nOnce the undesired portion is identified by the user, a filter script is automatically generated which removes the undesired Web page section. The filter script may be stored in a local file on the user\"\"s machine, or may be located remotely, for instance on an Internet service provider. The filter script may be associated with the URL of the Web page. Whenever the Web page is fetched by the user, any filter script which maps to the URL is automatically executed. Since some Web sites dynamically change a URL via CGI script arguments or load splitting across different servers, the URL associated with a filter script may be a regular expression. Thus, the filter may be applied to any URL which matches a URL pattern string. Actual filtering of the Web pages may be performed by a browser enabled with such abilities, or via a proxy process running on the browser\"\"s host."}
{"text": "The analysis of engine oil or other lubricant for the purpose of identifying premature component wearing has been performed for several decades using optical atomic spectroscopy (e.g., atomic emission spectroscopy (AES), as well as atomic absorption spectroscopy (AAS)). This technology was the basis for the military aviation's Spectroscopic Oil Analysis Program (SOAP). However, it has certain disadvantages, such as a lack of repeatability among different equipment and an inability to analyze particles greater than 5 μm in diameter. Furthermore, optical atomic spectroscopy is an elemental analysis of the total oil sample and typically does not characterize individual particles in the sample.\nOther approaches have since been proposed, whereby individual particles may be characterized and classified based on their chemical composition. However, these approaches are not sufficient to predict complex component failures."}
{"text": "According to conventional knowledge in the art of kraft pulping of cellulose, the level of dissolved organic materials (DOM) - - - which mainly comprise dissolved hemi-cellulose, and lignin, but also dissolved cellulose, extractives, and other materials extracted from wood by the cooking process - - - is known to have a detrimental affect in the later stages of the cooking process by impeding the delignification process due to consumption of active cooking chemical in the liquor before it can react with the residual or native lignin in wood. The effect of DOM concentration at other parts of cooking, besides the later stages, is according to conventional knowledge believed insignificant. The impeding action of DOM during the later stages of the cook is minimized in some state-of-the-art continuous cooking processes, particularly utilizing an EMCC.RTM. digester from Kamyr, Inc. of Glens Falls, N.Y., since the counter-current flow of liquor (including white liquor) at the end of the cook reduces the concentration of DOM both at the end of the \"bulk delignification\" phase, and throughout the so-called \"residual delignification\" phase.\nAccording to the present invention, it has been found that not only does DOM have an adverse affect on cooking at the end of the cooking phase, but that the presence of DOM adversely affects the strength of the pulp produced during any part of the cooking process, that is at the beginning, middle, or end of the bulk delignification stage. The mechanism by which DOM affects pulp fibers and thereby adversely affects pulp strength has not been positively identified, but it is hypothesized that it is due to a reduced mass transfer rate of alkali extractable organics through fiber walls induced by DOM surrounding the fibers, and differential extractability of crystalline regions in the fibers compared to amorphous regions (i.e. nodes). In any event, it has been demonstrated according to the invention that if the DOM level (concentration) is minimized throughout the cook, pulp strength is increased significantly.\nIt has been found, according to the present invention, that if the level of DOM is close to zero throughout a kraft cook, tear strength of the pulp is greatly increased, i.e. increased up to about 25% (e.g. 27%) at 11 km tensile compared to conventionally produced kraft pulp. Even reductions of the DOM level to one-half or one-quarter of their normal levels also significantly increase pulp strength.\nIn state-of-the-art kraft cooks, it is not unusual for the DOM concentration at some points during the kraft cook to be 130 grams per liter (g/l) or more, and at 100 g/l or more at numerous points during the kraft cook (for example in the bottom circulation, trim circulation, upper and main extractions and MC circulation in Kamyr, Inc. MCC.RTM. continuous digesters), even if the DOM level is maintained between about 30-90 g/l in the wash circulation (at later cook stages, according to conventional wisdom). In such conventional situations it is also not unusual for the lignin component of the DOM level to be over 60 g/l and in fact even over 100 g/l, and for the hemi-cellulose component of the DOM level to be well over 20 g/l. It is not known if the dissolved hemi-cellulose component has a stronger adverse affect on pulp strength (e.g. by adversely affecting mass transfer of organics out of the fibers) than lignin, or vice versa, or if the effect is synergistic, although the dissolved hemi-celluloses are suspected to have a significant influence.\nAccording to the present invention it has been recognized for the first time that the DOM concentration throughout a kraft cook should be minimized in order to positively affect bleachability of the pulp, reduce chemical consumption, and perhaps most significantly increase pulp strength. By minimizing DOM levels, one may be able to design smaller continuous digesters while obtaining the same throughput, and may be able to obtain some benefits of continuous digesters with batch systems. A number of these beneficial results can be anticipated by keeping the DOM concentration at 100 g/l or less throughout substantially the entire kraft cook (i.e., beginning, middle and end of bulk delignification), and preferably about 50 g/l or less (the closer to zero DOM one goes, the more positive the results). It is particularly desirable to keep the lignin component at 50 g/l or less (preferably about 25 g/l or less), and the hemi-cellulose level at 15 g/l or less (preferably about 10 g/l or less).\nAccording to the present invention it has also been found that it is possible to passivate the adverse affects on pulp strength of tile DOM concentration, at least to a large extent. According to this aspect of the invention it has been found that if black liquor is removed and subjected to pressure heat treatment according to U.S. Pat. No. 4,929,307 (the disclosure of which is hereby incorporated by reference herein), e.g. at a temperature of about 170.degree.-350.degree. C. (preferably 190.degree. C.) for about 5-90 minutes (preferably about 30-60 minutes) and then reintroduced, an increase in tear strength of up to about 15% can be effected. The mechanism by which passivation of the DOM by heat treatment occurs also is also not fully understood, but is consistent with the hypothesis described above, and its results are real and dramatic.\nAccording to the present invention-various methods are provided for increasing kraft pulp strength taking into account the adverse affects of DOM thereon, as set forth above, for both continuous and batch systems. Also according to the present invention increased strength kraft pulp is also provided, as well as apparatus for achieving the desired results according to the invention.\nAccording to one aspect of the present invention, a method of producing kraft pulp by cooking comminuted cellulosic fibrous material is provided. The method comprises the steps of continuously, at a plurality of different stages during kraft cooking of the material to produce pulp: (a) Extracting liquor containing a level of DOM substantial enough to adversely affect pulp strength. And, (b) replacing some or all of the extracted liquor with liquor containing a substantially lower effective DOM level than the extracted liquor, so as to positively affect pulp strength. Step (b) is typically practiced by replacing the withdrawn liquor with liquor selected from the group consisting essentially of water, substantially DOM free white liquor, pressure-heat treated black liquor, washer filtrate, cold blow filtrate, and combinations thereof. For example for at least one stage during cooking, black liquor may be withdrawn, and treated under pressure and temperature conditions (e.g. superatmospheric pressure at a temperature of about 170.degree.-350.degree. C. for about 5-90 minutes, and at least 20.degree. C. over the cooking temperature) to significantly passivate the adverse affects of DOM. The term \"effective DOM\" as used in the specification and claims means that portion of the DOM that affects pulp strength. A low effective DOM may be obtained by passivation, or by an originally low DOM concentration.\nThe method according to the invention can be practiced in a continuous vertical digester, in which case steps (a) and (b) are practiced at at least two different levels of the digester. There is also typically the further step (c) of heating the replacement liquor from step (b) to substantially the same temperature as the withdrawn liquor prior to the replacement liquor being introduced into contact with the material being cooked. Steps (a) and (b) can be practiced during impregnation, near the start of the cook, during the middle of the cook, and near the end of the cook, i.e. during substantially the entire bulk delignification stage.\nAccording to another aspect of the present invention a method of kraft cooking is provided comprising the steps of, near the beginning of the kraft cook: (a) Extracting liquor containing a level of DOM substantial enough to adversely affect pulp strength. And, (b) replacing some or all of the extracted liquor with liquor containing a substantially lower effective DOM level than the extracted liquor, so as to positively affect pulp strength.\nAccording to another aspect of the present invention a method of kraft cooking is provided comprising the steps of, during impregnation of cellulosic fibrous material: (a) Extracting liquor containing a level of DOM substantial enough to adversely affect pulp strength. And, (b) replacing some or all of the extracted liquor with liquor containing a substantially lower effective DOM level than the extracted liquor, so as to positively affect pulp strength.\nAccording to still another aspect of the present invention a method of kraft cooking pulp is provided comprising the following steps: (a) Extracting black liquor from contact with the pulp at a given cooking stage. (b) Pressure-heating the black liquor to a temperature sufficient to significantly passivate the adverse effects on pulp strength of DOM therein. And, (c) re-introducing the passivated-DOM black liquor back into contact with the pulp at the given stage.\nThe invention also comprises the kraft pulp produced by the methods set forth above. This kraft pulp is different than kraft pulps previously produced, having a tear strength as much as 25% greater at a specified tensile for fully refined pulp (e.g. at 9 km tensile, or at 11 km tensile) (and at least about 15% greater) compared to kraft pulp produced under identical conditions without the DOM maintenance or removal steps according to the invention, or as much as 15% greater (e.g. at least about 10% greater) where passified black liquor is utilized.\nThe invention is also applicable to kraft batch cooking of cellulosic fibrous material utilizing a vessel containing black liquor and a batch digester containing the material. In such a method of kraft batch cooking according to the invention there are the steps of: (a) Pressure-heating the black liquor in the vessel to a temperature sufficient to passivate the adverse effects on pulp strength of DOM therein. And, (b) feeding the black liquor to the digester to contact the cellulosic fibrous material therein. Step (a) is practiced to heat the black liquor at superatmospheric pressure at a temperature of about 170.degree.-350.degree. C. for about 5-90 minutes (typically at least about 190.degree. C. for about 30-60 minutes, and at least 20.degree. C. over cooking temperature), and step (b) may be practiced to simultaneously feed black liquor and white liquor to the digester to effect cooking of the cellulosic fibrous material.\nAccording to another aspect of the present invention an apparatus for kraft cooking cellulose pulp is provided. The apparatus comprises the following elements: An upright continuous digester. At least two withdrawal/extraction screens provided at different levels, and different cook stages, of the digester. A recirculation line and an extraction line associated with each of the screens. And, means for providing replacement liquor to the recirculation line to make up for the liquor extracted in the extraction line, for each of the recirculation lines. Each recirculatory loop typically includes a heater, and the digester may be associated with a separate impregnation vessel in which removal of high DOM concentration liquor and replacement with lower DOM concentration liquor also takes place (including in a return line communicating between the top of the impregnation vessel and the high pressure feeder).\nThe invention also relates to a commercial method of kraft cooking comminuted cellulose fibrous material by the step (a) of continuously passing substantially DOM-free cooking liquor into and out of contact with the material until completion of the kraft cook thereof, at a rate of at least 100 tons of pulp per day. This method is preferably practiced utilizing a batch digester, and by the further step (b), prior to step (a), of filling the digester with cellulose material, and the further step (c), after step (a) of discharging kraft pulp from the digester. The invention also relates to a batch digester system for practicing this aspect of the invention.\nIt is the primary object of the invention to produce increased strength kraft pulp. This and other objects of the invention will become clear from an inspection of the detailed description of the invention and from the appended claims."}
{"text": "As an apparatus for driving a rotary machine such as a compressor, there is an apparatus including an electric driving device for generating a rotational driving force and a transmission device for changing the speed of a rotational driving force generated by the electric driving device and then transmitting the speed-changed rotational driving force to the rotary machine.\nPatent Document 1 discloses that a constant-speed motor and a variable-speed motor for speed change are used as the electric driving device and a planetary gear transmission device is used as the transmission device to accurately control a gear ratio."}
{"text": "The present invention relates generally to power switching devices. More particularly, the present invention relates to a current sensing solution for power switching devices which reduces power losses and parasitic inductance without correspondingly occupying excessive space in a power supply.\nHigh efficiency is one of the important parameters in designs for modern power supplies, and power density has increased accordingly. A feedback (sensed) current signal is another important factor in proper operation of the power conversion circuit.\nOne conventional example is represented in FIG. 1 of a current sensing circuit 10 for a power conversion circuit implementing MOSFET devices 11 (having gate 12, drain 13 and source terminals 14), wherein the current signal is obtained from opposing ends 15, 16 of a current sensing resistor.\nAnother conventional example of a current sensing circuit 20 as represented in FIGS. 2a and 2b may include a discrete current sense transformer 21 having a primary winding coupled to the drain terminal of a MOSFET switching device 11 and a secondary winding to which current sensing leads 15, 16 are coupled.\nOne of skill in the art may appreciate that each of the two aforementioned examples contribute to additional losses due to the copper winding/track and soldering joint, and also introduction of parasitic inductance due to additional copper track. Such losses arise in the context of power switching devices implementing either fully controllable power switching (e.g., three terminal devices such as a MOSFET), or uncontrollable power switching (e.g., two terminal devices such as a power diode).\nYet another problem which arises in such conventional arrangements is undesirable occupation of physical space, particularly with respect to the discrete transformer, whether with respect, for example, to horizontal surface area on a printed circuit board or vertical cross-section in a stacked arrangement."}
{"text": "The development of a method for decomposing or removing nitrogen oxides or other harmful pollutants in the exhaust gas from automotive vehicles has been increasingly demanded due to environmental pollution. A three way catalytic converter, for removing carbon monoxide, hydrocarbons and nitrogen oxides concurrently, is a practical catalyst for purification of the exhaust gas. Cordierite coated with gamma-alumina or other heat resistant carrier supporting Pd, Pt, Rh or other noble metal is for general use.\nHowever, in a lean-burn engine and a diesel engine, which are operated under the condition of high oxygen concentration the three way catalytic converter is ineffective because the noble metals are poisoned with excess oxygen. Although zeolites have been studied as the catalyst serving under the condition of high oxygen concentration, their structure easily breaks under the hydrothermal conditions. Oxide catalysts superior in heat resistance have also been studied, however, it is insufficient in catalytic activity for their practical use.\nOn the other hand, in thermal power plants or other fixed combustion equipments, nitrogen oxides in the exhaust gas containing excess oxygen are reduced by the oxide catalyst with ammonia as a reductant. Because ammonia has toxicity, a safety problem is created when this system is utilized for automotive vehicles or other moving combustion equipments or the fixed combustion equipments installed in an urban area.\nRecently, a method of removing nitrogen oxides in the exhaust gas containing excess oxygen was developed. In the method, nitrogen oxides are absorbed by an absorbent in the lean-burn region or in the presence of excess oxygen, and nitrogen oxides are desorbed and reduced by means of the three way catalyst in the rich-burn region or in the vicinity of the stoichiometric air/fuel ratio.\nKnown absorbents of nitrogen oxides are complex oxides of the Ba--Cu--O system and the Mn--Zr--O system. However, a material having better and stable absorbability even at high temperatures is demanded."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus and method for sensing undesired flammable vapors at or near gaseous fuel-fired burners such as may be used for gas-fired appliances, for example residential gas-fired water heaters.\n2. Description of Prior Art\nSafety issues are important considerations in the design, manufacture and operation of gaseous fuel-fired burners, such as found in gas-fired appliances. This is especially true for appliances operated in a residential environment. Such appliances include residential gas-fired water heaters, gas-fired clothes dryers, furnaces, etc.\nOne safety issue involves a potential danger which arises when burners are exposed to flammable vapors in the ambient atmosphere external to the appliance. Such flammable vapors can cause uncontrolled propagation of flames and/or possible explosions.\nAn object of this invention is to provide a burner for a gas-fired appliance, capable of detecting and responding to the presence of undesired flammable vapors at, near or in the vicinity of the burner. One goal would be to turn down and/or completely shut off a gaseous fuel supply to the burner, to prevent an undesired propagation of flame and/or an explosion.\nAnother object of this invention is to provide a flammable vapor sensing burner which has a high degree of sensitivity to respond to even very low concentrations of detected flammable vapor.\nStill another object of this invention is to provide a burner having a simplified construction with low cost, long life expectancy and enhanced durability of construction and reliability of operation.\nThese and other objects of this invention will become apparent in light of the present specification, the claims and the drawings.\nThe apparatus of this invention includes a pilot burner that operates with a main burner of a gas-fired appliance, for sensing flammable vapor present in an atmosphere ambient to the gas-fired appliance.\nA first fuel gas transporting member communicates with and is connectable to a source of fuel gas. As used throughout this specification and in the claims, the phrase in communication with is intended to relate to two spaces or voids fluidically connected with or communicating with each other. The source of fuel gas may be configured for providing fuel gas, through a second fuel gas transporting member, to a main burner of the gas-fired appliance.\nAt least one mixer member preferably mixes fuel gas received by the first fuel gas transporting member, with air taken from the atmosphere external and ambient to the appliance, which air may contain an undesired level or amount of flammable vapor. A fuel gas-air transporting member may be operably connected to the mixer member or another suitable means for mixing fuel gas. The fuel gas-air transporting member preferably is configured for receiving the mixed fuel gas and air which may contain flammable vapor, and transporting the mixed fuel gas and air to a combustion chamber within the gas-fired appliance.\nA primary pilot flame holder may be operably positionable within the combustion chamber of a gas-fired appliance. The primary pilot flame holder is preferably in communication with or operably connected to the fuel gas-air transporting member.\nAt least one secondary air shield preferably operates with the primary pilot flame holder, for substantially precluding access of secondary air which may be present within the combustion chamber of a gas-fired appliance, to a flame established on the primary pilot flame holder.\nAn igniter or another suitable ignition means is positioned near and operably associated with the at least one secondary air shield, for facilitating ignition of a main burner of a gas-fired appliance.\nAt least one sensor detects a change in a characteristic of a primary pilot flame established on or held by the primary pilot flame holder, such change resulting from a presence of and/or a change in an amount of, flammable vapor in the air of the atmosphere ambient to the gas-fired appliance.\nThe at least one sensor is sufficiently sensitive to detect a change corresponding to preselected relative amount of flammable vapor present in the air of the atmosphere ambient to the gas-fired appliance.\nA controller preferably is in electrical communication with or operably associated with the at least one sensor for detecting a change in a characteristic of a flame, for controlling the operation of the source of fuel gas for a gas-fired appliance, in response to the presence of and/or the change in the amount of flammable vapor in the air of the atmosphere ambient to the gas-fired appliance. Upon demand, the controller can cease or halt the supply of fuel gas to the gas-fired appliance, for example when the at least one sensor detects a change corresponding to a threshold amount or a predetermined amount of flammable vapor in the air ambient to the gas-fired appliance.\nIn a preferred embodiment of this invention, the at least one mixer member for mixing fuel gas with air taken from the atmosphere external and ambient to the gas-fired appliance, which air may contain undesired flammable vapor, comprises a mixer tube positioned substantially adjacent to but spaced apart from an outlet of the first fuel gas transporting member. An inlet of the mixer tube is preferably positioned to receive a fuel gas jet from the gas nozzle. The outlet of the first fuel gas transporting member and the mixer tube are preferably operably positioned substantially external to a gas-fired appliance, so that air from the atmosphere ambient to the gas-fired appliance is entrained with the fuel gas jet from the outlet of the first fuel gas transporting member, and transported into the inlet of the mixer tube.\nThe at least one secondary air shield may comprise a secondary air shield, operably positioned substantially about the primary pilot flame holder, so as to substantially preclude access by secondary air within the combustion chamber of the gas-fired appliance to a pilot flame held by the primary pilot flame holder, so that any flame is substantially unaffected by any such secondary air.\nThe at least one sensor for detecting a change in a characteristic of a flame held by the primary pilot flame holder may comprise a thermocouple operably disposed proximate to the primary pilot flame holder for sensing the temperature at, near or in the immediate vicinity of the primary flame holder, and for providing a signal representative of the temperature which has been sensed.\nIn one embodiment of this invention, the controller which operates the source of fuel, in response to the presence of flammable vapor in the air ambient to the gas-fired appliance may comprise a gas valve connected directly to the thermocouple, so that as the thermocouple cools the gas valve closes.\nIn an alternative embodiment of this invention, the controller for controlling the operation of the source of fuel, in response to the presence of flammable vapor in the air ambient to the gas-fired appliance may comprise means for comparing the signal representative of the temperature in the immediate vicinity of the primary pilot flame holder. The means for comparing may be operably configured to compare the signal provided by the thermocouple to a predetermined range of values representative of acceptable conditions for the flame established on or held by the primary pilot flame holder. The source of fuel gas can be controlled to cease or halt the flow of fuel gas to the main burner and pilot burner when the signal provided by the thermocouple departs from a predetermined range of values representative of or corresponding to acceptable conditions for the flame held by the primary pilot flame holder.\nIn an alternative embodiment of this invention, in response to the presence of flammable vapor in the air ambient to the gas-fired appliance the controller may further comprise means for determining the rate of change over time of the temperature at, near or in the immediate vicinity of the primary flame holder. The determined rate of change over time of the temperature can be compared to a preselected value for the rate of change. Means may also be provided for reducing or stopping flow of the fuel gas to the main burner and/or the pilot burner when the rate of change over time of the temperature is both negative and exceeds a predetermined absolute value.\nIn a preferred embodiment of this invention, the at least one secondary air shield is a vertically oriented, substantially hollow cylindrical member having a primary pilot flame holder disposed at a lower end and a substantially open upper end opposite the lower end. The at least one secondary air shield can be positioned proximate the main burner flame holder of the gas-fired appliance. The means, operably associated with the at least one secondary air shield, for igniting the main burner of the gas-fired appliance comprises a secondary pilot flame holder surface, operably disposed about the substantially open upper end of the substantially hollow cylindrical member, for establishing a secondary pilot flame and for igniting fuel gas migrating from a main burner flame holder, toward establishing a main burner flame.\nIn another preferred embodiment of this invention, the at least one secondary air shield is a horizontally oriented, substantially hollow cylindrical member, having a primary pilot flame holder disposed at one end and a substantially closed end opposite the one end. The at least one secondary air shield can be positioned proximate the main burner flame holder of the gas-fired appliance. The means, operably associated with the at least one secondary air shield, for igniting the main burner of the gas-fired appliance may comprise a secondary pilot flame holder aperture, operably disposed in an upwardly facing region of a cylindrical side portion of the substantially hollow cylindrical member, for igniting fuel gas migrating from a main burner flame holder, toward establishing a main burner flame.\nIn still another preferred embodiment of this invention, the at least one secondary air shield is a substantially hollow, cylindrical member, having a primary pilot flame holder positioned at one end. The at least one secondary air shield can be positioned proximate the main burner flame holder of the gas-fired appliance. The means, operably associated with the at least one secondary air shield member, for igniting the main burner of the gas-fired appliance comprises at least one flash tube member, operably associated with or in communication with the substantially hollow cylindrical member, for receiving fuel gas from a main burner flame holder and for conducting the fuel gas toward a primary pilot flame held by the primary pilot flame holder, toward igniting the fuel gas from the main burner flame holder and thus establishing the main burner flame.\nThis invention also relates to a method for sensing flammable vapor present in the atmosphere ambient to the gas-fired appliance, and for controlling operation of the gas-fired appliance in response to a presence of or an undesired level of a flammable vapor.\nA first fuel gas transporting member is connected to or communicates with a source of fuel gas. The source of fuel gas can be configured to provide fuel gas through a second fuel gas transporting member, to the main burner of the gas-fired appliance. Fuel gas received by the first fuel gas transporting member is mixed, using least one mixer member, with air taken from the atmosphere external and ambient to the appliance, which air may contain flammable vapor.\nThe mixed fuel gas and air is transported or flows to a combustion chamber within the gas-fired appliance, such as through a fuel gas-air transporting member. A primary pilot flame holder operably positionable within the combustion chamber of the gas-fired appliance is connected to or communicates with the fuel gas-air transporting member.\nAt least one secondary air shield is used to substantially preclude access of secondary air which may be present within the combustion chamber of the gas-fired appliance to the flame established on or held by the primary pilot flame holder.\nAt least one sensor is used to detect a change in a characteristic of a flame held by the primary pilot flame holder, such change resulting from the presence of and/or the change in the amount of flammable vapor in the air from the atmosphere ambient to the gas-fired appliance. The change corresponds to a preselected relative amount of flammable vapor present in the air from the atmosphere ambient to the gas-fired appliance.\nA controller is used to control the operation of the source of fuel gas for a gas-fired appliance, in response to the presence of and/or the change in the amount of flammable vapor in the air from the atmosphere ambient to the gas-fired appliance. The supply of the fuel gas to the gas-fired appliance is ceased or halted when the at least one sensor detects a change corresponding to a predetermined amount of flammable vapor in the air ambient to the gas-fired appliance.\nIn one preferred embodiment, a mixer tube is positioned substantially adjacent to but also spaced apart from an outlet of the first fuel gas transporting member. An inlet of the mixer tube can be positioned to receive the fuel gas jet from the gas nozzle.\nThe outlet of the first fuel gas transporting member and the mixer tube is preferably positioned substantially external to the gas-fired appliance, so that air from the atmosphere ambient to the gas-fired appliance is entrained with the fuel gas jet from the outlet of the first fuel gas transporting member, and directed or transported into the inlet of the mixer tube.\nA secondary air shield can be positioned substantially about the primary pilot flame holder, so as to substantially preclude access of secondary air within the combustion chamber of the gas-fired appliance to the pilot flame established on or held by the primary pilot flame holder, so that the pilot flame is substantially unaffected by any such secondary air.\nA thermocouple can be positioned proximate to the primary pilot flame holder for sensing the temperature at, near or in the immediate vicinity of the primary pilot flame holder, and for providing a signal representative of the temperature of the immediate vicinity of the primary pilot flame holder.\nThe thermocouple can be connected directly to a gas valve, so that as the thermocouple cools and reaches a certain level the gas valve then closes.\nThe signal representative of the temperature at, near or in the immediate vicinity of the primary pilot flame holder is preferably compared to a predetermined range of values representative of or corresponding to acceptable conditions for the flame held by the primary pilot flame holder.\nThe source of the fuel gas ceases or halts fuel gas flow to the main burner and/or the pilot burner when the signal provided by the thermocouple departs from or is different than the predetermined range of values.\nThe rate of change overtime of the temperature at, near or in the immediate vicinity of the primary flame holder is determined and compared to a preselected value for the rate of change over time, of the temperature.\nThe source of the fuel gas then ceases or halts fuel gas flow to the main burner and/or the pilot burner preferably when the rate of change over time of the temperature is both negative and exceeds a predetermined absolute value."}
{"text": "The present invention relates to a motor-generator set and more particularly to a three-phase constant frequency motor-generator having only one single rotor for both the motor and the generator, the output frequency of the generator being constant notwithstanding the fact that the rotational speed of the rotor may vary."}
{"text": "1. Field of the Invention\nThe present invention relates generally to high-frequency mixers and more specifically, to an odd-subharmonic balanced mixer for sub-millimeter-waves. This new mixer yields improved conversion efficiency, lower output noise, superior cross-port isolation, improved impedance matching, and ease of fabrication at receiver frequencies on the order of several hundred gigahertz.\n2. Prior Art\nA mixer is a critical component of modern radio-frequency (RF) systems. It is usually the first or second device after the RF antenna input. Various mixer parameters, such as bandwidth and inter-port isolation, must be optimized to produce devices capable of performing in modern RF systems.\nIn recent years there has been an increase in sub-millimeter-wave receiver applications requiring easily producible and improved performance mixer technology at these extremely high frequencies.\nOne reason frequency conversion is done is to make signal processing easier and less expensive. Changing the frequency of a signal without altering its information content is necessary because signal processing components, such as amplifiers, are much less expensive and work better when they are designed to operate at lower frequencies.\nAnother reason for frequency conversion is to allow for the practical transmission of low-frequency information (such as audio information) through space. Transmitting audio frequency (up to 20 KHz) signals directly would require antennas of impractically large size because of the relatively long wavelengths of audio signals. However, if the audio signals are first converted up in frequency to a much higher carrier frequency, antennas of practical size can be utilized. At the receiving end, it is necessary to capture part of the electromagnetic energy of the transmitted signal and reconvert it back down to the audio frequency range in order to extract the original low-frequency information. Thus, both transmission and receiving require an input signal to be converted. Mixers perform this frequency conversion.\nMixing an input RF signal with an LO signal yields frequency products below and above the RF and LO frequencies. Each frequency product corresponds to the sum of the input RF and LO frequencies, while the lower frequency product corresponds to the difference between the input RF and LO frequencies. Usually, it is the lower frequency product (the xe2x80x9cdown-convertedxe2x80x9d RF signal) which is used in receiving systems, whereas the high frequency product (the xe2x80x9cup-convertedxe2x80x9d RF signal) which is used in transmission systems. A non-linear element, e.g., diode, is essential for frequency conversion of any sort.\nThere are basically four types of mixers: each with a different number of diodes: single-ended, single-balanced, double-balanced, and double double-balanced (also called triple-balanced). However, all types are three-port devices and comprise an input port (the RF port), a local oscillator input port (the LO port), and an output port (the IF port). Single-ended mixers are the simplest type and are realized using only a single diode. The LO, RF and IF ports are separated only by filters to provide some degree of inter-port isolation. Single-ended mixers, however, have a narrow bandwidth, limited dynamic range and poor inter-port isolation.\nBroader bandwidths and better isolation can be obtained with a single-balanced mixer. A single-balanced mixer consists of two single ended mixers. The mixer diodes are fed by the LO and RF signals through a BALUN which yields inter-port isolation between the LO and RF ports. Double-balanced mixers are so called because they have two BALUNs instead of just one and comprise two single-balanced mixers. This enables inter-port isolation both between the LO and RF ports and between the LO and IF ports. Double-balanced mixers use twice the number of diodes (four) as a singly-balanced mixer and the diodes are usually, although not always, arranged in a diode xe2x80x9cquadxe2x80x9d ring configuration.\nTriple-balanced mixers are so called because, in addition to BALUNs on the RF and LO ports, the IF port is also balanced. A triple-balanced mixer requires twice the number of diodes (eight) as a doubly-balanced mixer and a triple-balanced mixer is usually realized by combining two diode quad ring mixers.\nThe mechanism by which a diode mixer converts energy from one frequency to another is generally well known. A diode which is pumped by a signal from a local oscillator at a frequency fLO responds as a switch that is closed during the conducting portion of the LO voltage cycle and open during the non-conducting part of the cycle. This switching action modulates an incident RF signal at a frequency fRF and thereby makes available output signals at numerous frequencies different from and in addition to the incident RF and LO signals. The relationship between the various additional output signals is fIF=xc2xdkfLOxc2x1fRFxc2xd, where fIF is the output frequency and k and I are integers (at low level external inputs I=1). If, for example, k=1, the mixer is considered a fundamental mixer.\nHarmonic mixers have been utilized in which k greater than 1, hence the principal output is fIF=xc2xdmfLOxc2x1fRFxc2xd that is, at a particular xe2x80x9cmxe2x80x9d th harmonic of the LO frequency. These mixers have used filters to block responses corresponding to k 1m. The main disadvantage to this type of mixer is that the conversion-loss is relatively high because of the limited success of filters to prevent RF signal energy to convert into unintended outputs.\nA more recent type of even subharmonic (2n) pumped mixer uses two diodes connected in parallel and opposing polarity, that is, antiparallel (n is any integer). These form a semiconductor switch which is toggled at 2nfLO and, consequently, output frequencies fIF are available only as combination of even harmonics of the local oscillator frequency.\nAs background, existing forms of the two-diode, subharmonic pumped mixers are described in literatures such as:\nE. R. Carlson et al, xe2x80x9cSubharmonically Pumped Millimeter-Wave Mixers,xe2x80x9d Vol. MTT-26, pp. 706-715, October 1978;\nM. V. Schneider and W. W. Snell, Jr., xe2x80x9cHarmonically Pumped Stripline Down-Converter,xe2x80x9d Vol. MTT-23, pp. 271-275, March 1975; and\nM. Cohn et al, xe2x80x9cHarmonic Mixing With An Antiparallel Diode Pair,xe2x80x9d Vol. MTT-23, pp. 667-673, August 1975.\nSaid xe2x80x9cMTTxe2x80x9d volumes being the IEEE Transactions on Microwave Theory and Techniques.\nAn example of a balanced mixer utilizing a subharmonic pumped antiparallel diode pair is disclosed in U.S. Pat. No. 3,983,489, issued Sept. 28, 1976 to Gittinger. This type of mixer, in which transformer coupling is used to isolate the RF and LO ports, can achieve wide bandwidth in the UHF frequency range but is not usable at microwave frequencies due to limitations of the transformer devices used therein.\nA search of the prior art has revealed the following U.S. patents which appear pertinent in varying degrees to the present invention:\nU.S. Pat. No. 3,311,811 Rupp\nU.S. Pat. No. 4,145,692 Armstrong et al\nU.S. Pat. No. 4,320,536 Dietrich\nU.S. Pat. No. 4,392,255 Del Giudice\nU.S. Pat. No. 4,394,632 Hu\nU.S. Pat. No. 4,485,488 Houdart\nU.S. Pat. No. 4,491,977 Paul\nU.S. Pat. No. 4,730,169 Li\nU.S. Pat. No. 4,803,419 Roos\nU.S. Pat. No. 4,817,200 Tanbakuchi\nU.S. Pat. No. 4,994,755 Titus et al\nU.S. Pat. No. 5,060,298 Waugh et al\nU.S. Pat. No. 5,077,546 Carfi et al\nU.S. Pat. No. 5,266,963 Carter\nU.S. Pat. No. 5,446,923 Martinson et al\nU.S. Pat. No. 5,517,687 Niehenke et al\nU.S. Pat. No. 5,553,319 Tanbakuchi\nU.S. Pat. No. 5,721,514 Crockett et al\nU.S. Pat. No. 5,740,528 Drennen\nU.S. Pat. No. 5,819,169 Fudem\nU.S. Pat. No. 5,844,939 Scherer et al\nOf the aforementioned patents, the following appear to be the most relevant:\nU.S. Pat. No. 5,553,319 to Tanbakuchi is directed to a routing YIG tuned mixer. A routing YIG tuned resonating filter integrated with an image enhanced double-balanced mixer is provided. FIG. 10 is a schematic diagram of an image enhanced double-balanced YIG tuned mixer 305 which mixes a fundamental or odd-harmonic products of the local oscillator signal with the RF input signal to produce a predetermined low frequency IF signal output. As shown in FIG. 10, the mixer 305 comprises a second half loop of coupling loop 138 and a full output coupling loop 142 in which coupling loop 142 is continuous at a first end and connected to low output 194. Output coupling loop 142 is terminated at a second end by a monolithic GaAs diode integrated circuit comprising a Schottky diode bridge. Mixer 305 further comprises a single IF BALUN 306. IF BALUN 306 functions as a fundamental and odd harmonic IF BALUN.\nU.S. Pat. No. 5,446,923 to Martinson et al is directed to a mixer using fundamental frequency or second or third harmonic frequencies of a local oscillator for maximized resultant frequency mixer product. This device shows that the third harmonic mixing can be accomplished provided that the mixing element is a two terminal device having a substantially symmetrical non-linear forward and reverse voltage current characteristic. Also required is that it have a pair of conduction threshold voltages which are substantially symmetrical above and below zero volts, beyond which the mixing element will be conductive at least when a signal from a local oscillator is imposed on the mixing element. The idealized reflection coefficient waveform has a fundamental frequency that is equal to the fundamental frequency of the local oscillator when the second DC component of bias voltage is imposed across the mixing element so as to give the maximized resultant frequency fINxc2x1fLO. The idealized reflection coefficient waveform has a large amplitude component that is at a frequency which is three times the fundamental frequency of the local oscillator when the third DC component of bias voltage is imposed across the mixing element so as to give the maximized resultant frequency fINxc2x1fLO.\nU.S. Pat. No. 4,994,755 To Titus et al is directed to an active BALUN. FIG. 1 shows an active double BALUN that may be used in a double BALUN mixer. With such an arrangement, the double BALUN provides two outputs, RFLO and RFLO. These signals are down-converted to an IF frequency, the difference frequency between RF and LO. Any unconverted RF or LO will substantially cancel at the output of the mixer because of the 180 degree phase shift difference. Even order harmonic products will be substantially canceled by this arrangement. The harmonic generation occurs in the mixer according to the formula IF=xc2x1nRFxc2x1mLO. The harmonic order for a doubly balanced mixer refers to the harmonic order product n+m. For example, 2RF-LO is a third order odd harmonic product. The fourth order even harmonic product is 2RF-2LO. The even order product is canceled by this double balanced mixer arrangement. FIG. 11 shows an embodiment of this device in which active devices at the single input BALUN operate in the non-linear mode to produce harmonics of the RF. With this arrangement, even harmonics of the RF will be canceled and odd harmonics of the RF will be coupled to the output terminal.\nU.S. Pat. No. 4,817,200 to Tanbakuchi is directed to a tracking YIG tuned filter mixer. FIG. 5 shows the YIG resonator 304 and its associated circuitry function to mix the filtered RF signals which appear on line 305 with the local oscillator signals on line 319. The input of full loop 381 is connected to line 319 which carries the local oscillator signal. Two diodes 385 and 387 are connected in series across the output side of full loop 381 and output line 317 is connected to the junction between the diodes. A matching network 391 is connected to line 317 near the junction point between the diodes. This mixer is an odd harmonic mixer in which the IF signal at a frequency of FIF=(2+1) (FLO)xc2x1FRF flows through the local oscillator IN port and the IF output port.\nU.S. Pat. No. 5,844,939 to Scherer et al is directed to a low-cost phase-locked local-oscillator for millimeter-wave transceivers. The transceiver includes a transceiver oscillator and a receiver mixer. The transceiver oscillator generates a transceiver drive signal. The receiver\"\"s mixer mixes a received modulated signal with a harmonic of the transceiver drive signal to generate an intermediate frequency receive signal. FIG. 2 shows an embodiment in which the transceiver drive mixer 44 is a harmonic mixer in which an odd harmonic of the intermediate transmit oscillator signal is mixed with the transceiver drive signal. Typically, the third or fifth harmonic of the intermediate transmit signal is mixed with the transceiver signal.\nThe invention accomplishes a high performance mixing function with a third harmonic of a local oscillator as part of a heterodyne receiver-facilitating the utilization of lower frequency local-oscillators for higher frequency receivers (e.g., a 150 GHz local-oscillator in a 450 GHz receiver). The invention generally overcomes the technological gap created between rapid increase in sub-millimeter-wave receiver applications, and slow advancements in local-oscillator technology.\nFurthermore, specifically, this invention also overcomes mixer fabrication difficulties associated with implementations of balanced mixers at millimeter and sub-millimeter wave frequenciesxe2x80x94demonstrating a simple implementation of an all waveguide solution.\nThe inventive balanced sub-harmonic mixer features performance advantages relative to single-ended sub-harmonic mixers operating with an identical frequency harmonic of the local-oscillator. The two-diode balanced mixer configuration features symmetry which suppresses, without filters, a significant amount of undesired frequency down conversion combinations in favor of desired ones, yielding better conversion-loss, lower output noise, and superior cross-port isolation, relative to one diode conventional sub-harmonic mixers. The two-mixer diode balanced mixer configuration also yields significantly easier design requirements at the input and output mixer ports yielding easier, wider and better impedance matching relative to one diode conventional sub-harmonic mixers.\nThis invention embodies a unique, easy to produce implementation of a balanced mixer sub-harmonically driven, particularly attractive for millimeter-wave and sub-millimeter-wave receivers.\nThe relationships between E-vectors and mixer diodes (mounted on an additional printed-circuit on dielectric substrate at the BALUN) are at the xe2x80x9cheartxe2x80x9d of this mixer invention. These relations facilitate diodes\"\" LO drives out-of phase, and the series distribution of the RF signal on the same pair of diodes.\nThis invention embodies a unique implementation of a balanced mixer circuit structure facilitating ease of fabrication of balanced mixers at millimeter and sub-millimeter wave frequencies. Particularly, this invention embodies a balanced mixer implemented as an odd-sub-harmonic balanced mixerxe2x80x94facilitating xe2x80x9cstretchingxe2x80x9d useful employment of lower frequency LOs for higher band receivers. At the xe2x80x9cheartxe2x80x9d of the invention is a unique realization of a balanced-unbalanced (BALUN) networkxe2x80x94a network required for implementing balanced mixers.\nIt is a principal object of the present invention to provide a balanced mixer for sub-millimeter frequencies practical to produce and which yields lower conversion loss, lower noise, better cross-port isolation and impedance matching.\nIt is another object of the invention to provide a two-mixer diode balanced mixer using an augmented sub-harmonic driven local oscillator for high frequency receivers.\nIt is still an additional object of the invention to provide an odd-sub-harmonic balanced mixer implemented with a mixer diode assembly positioned within a mechanical xe2x80x9callxe2x80x9d waveguide structure."}
{"text": "Data networks are used to deliver content and services to end users, as well as to internally transport data within a group of networks. The route between a content provider or service provider and a destination of a given packet flow often can involve multiple networks (also known as autonomous systems). The operators of such networks typically have agreements with each other so as to arrange for the transport of packet flows between an originating network and a destination via one or more intermediary networks. Typically these agreements are “settlement free” in that the data is transferred between networks without requiring financial compensation for the data traffic transmitted by one network on behalf of another network.\nThe use of the same reference symbols in different drawings indicates similar or identical items."}
{"text": "It is known to use sensors coupled to a mattress for monitoring a bed occupant. For example, Triplett et al. (U.S. Pat. No. 4,175,263) disclose a pressure sensing mechanism for placement on a hospital bed. The pressure sensing mechanism comprises a first pressure sensing pad for sensing a patient's weight when the patient is near the centre of the bed, and a second pressure sensing pad for sensing the patient's weight when the patient is near the edges of the bed. An alarm may be triggered when the patient moves to a location near the edge of the bed.\nTucknott et al. (U.S. Pat. No. 4,633,237) disclose a patient bed alarm system comprising a matrix of sensors woven into a mat for placement on a bed. The matrix of sensors is coupled to a micro computer for calculating the position of a patient in the bed. An alarm may be triggered when the patient is about to leave the bed.\nAlihanka et al. (U.S. Pat. No. 4,320,766) disclose a capacitive motion sensor placed under a mattress or the like for monitoring the movements of a person.\nMcMahon et al. (U.S. Pat. No. 5,435,317) disclose a device for detecting a respiratory dysfunction of a person located in a bed. The device comprises a detection unit which is provided under the mattress of the bed, and a stimulation unit, which acts, in response to a signal from the detection unit, to impart a rocking motion to the bed. The detection unit is embodied as a pad-like device provided under the mattress.\nBellin et al. (U.S. Pat. No. 5,558,996) disclose a patient monitoring device wherein sensors are located in a bed sheet with which a subject comes in contact. One sensor produces a signal corresponding to respiratory induced, pulmonary motion, and myocardial pumping sounds. A second sensor produces a signal corresponding to changes in body position. A processor amplifies and filters the induced signals resulting in resolved output highly correlated to respiration rate, heart beat rate, and changes in body position.\nReimer et al. (U.S. Pat. No. 5,917,180) disclose a pressure sensor which relies on detecting multiply scattered light within an optical cavity. Changes in the volume of the cavity are sensed by the change in sampled light intensity. Pressure sensitive mats with a high density of sensors are assembled using optical fiber technology.\nMusick (U.S. Pat. No. 5,844,488) discloses a narrow, pressure-sensitive sensor pad for installation on top of and across the width of a mattress proximate the midsection of a reclining patient. The pad has both central and edge switching areas. The central pressure sensitive switch indicates the presence of a patient in the center of the bed. When a patient moves toward either edge of the bed, an edge switch is activated which generates an early warning signal indicating to attending personnel that a patient has moved from the center of the bed to an edge and may be attempting to exit the bed unattended.\nHammett (U.S. Pat. No. 5,144,284) discloses a bed covering device adapted to cover a mattress. The device has compliant flat pressure sensitive means disposed its lower surface for detecting compressive force exerted by a person, and electrical connector means emergent from the pressure sensitive means and configured to connect in modular fashion with a monitoring device and alarm.\nRudeke (PCT publication No. WO9010281A1) discloses a bed alarm intended for use with patients in hospitals and like establishments for indicating when a patient leaves his or her bed. The alarm includes a pressure sensor which functions to detect whether a patient is lying in bed or not and which is intended to be positioned in the bed and preferably between the bed-bottom and the mattress.\nRincoe et al. (U.S. Pat. No. 5,993,400) and Scott (U.S. Pat. No. 6,067,019) disclose devices that utilize arrangements of sensors to detect the impending egress of a bed occupant. Joseph et al. (U.S. Pat. No. 5,410,297) disclose a weight-sensitive capacitive sensor for tracking the position of an occupant. Dixon et al. (U.S. Pat. No. 6,208,250) disclose a system comprising two sensors and a processor used to determine the location of a patient on a bed.\nRosenthal (U.S. Pat. No. 3,961,201) discloses a switch located between a mattress and a bedframe that is used to detect when a patient moves close to the edge of a bed. Feldl (U.S. Pat. No. 4,020,482) discloses an air bladder below a mattress that signals egress of a patient when air pressure in the bladder falls below a threshold. Nicholas (U.S. Pat. No. 4,381,434) discloses a device including a spring-loaded plate with a limit switch, all mounted below a mattress.\nDespite the volume of existing patent literature, there are significant shortcomings in the known devices. For many devices, reliable operation requires installation on top of the mattress. This has two major shortcomings: firstly, the devices are often stiff (i.e. resistant to flexing) and hard, thereby causing discomfort to the bed occupant. In some applications (for example, monitoring of demented occupants) this can cause physical and psychological irritation to the bed occupant. Secondly, because the devices are on top of the mattress, they rapidly suffer material fatigue and fail under normal use."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrical connector, more particularly to an electrical connector with a modular design.\n2. Description of the Prior Art\nMost of PCs (personal computer), TV sets and electronic devices have an outer casing provided with built-in electrical connector for electrical connection with a peripheral device (such as DVD player) to facilitate signal transmission therebetween.\nReferring to FIG. 1, a conventional electrical connector 100a is shown to include a metal shell 130a, an insulated main body 110a, and a plurality of terminals 120a. The insulated main body 110a is formed with a plurality of retention holes 111a. The terminals 120a are inserted respectively into the retention holes 111a in the insulated main body 110a. Afterward, the insulated main body 110a is inserted into the metal shell 130a, thereby completing assembling of the conventional electrical connector 100a. \nIn case the terminals 120a are not aligned with one another along a horizontal plane at an initial condition (i.e. prior to inserting into the main body 11a), collision of the terminals 120 against the peripheral walls defining the retention holes 111a respectively during the insertion may result in partial bending of the terminals and hence disqualified product after assembly, which must be discarded, thereby occurring an extra manufacture expense.\nMoreover, the metal shell 130a is usually fabricated according to different designs and different specifications, and has a structure differ from one another. In case, the dimension or design of the metal shell 130a is abruptly changed, the finished insulated main body 110a is unable to complement with the metal shell 130a. The finished insulated main body 110a and the metal shell 130a must be discarded. To open a new mold for fabricating a required insulated main body will cause extra expense for the manufacturers."}
{"text": "As methods of manufacturing such plant cultivating substrate, there are known a method of manufacturing it by mixing and reacting water-retentive filling material such as bark compost, chaff compost with aqueous solution or solvent dispersion of a small amount of urethane prepolymer and then allowing the mixture to cure, and a method of manufacturing by reacting and curing a suspension containing water-retentive filling material, water and urethane prepolymer (see, e.g. Patent Documents 1 and 2).\nThere is also known a method by pre-mixing water-retentive filling material and urethane prepolymer and then adding a small amount water to the resultant mixture for its reaction and curing (see e.g. Patent Document 3).    Patent Document 1: Japanese Patent Application Published No. 56-18165    Patent Document 2: Japanese Patent Application “Kokai” No. 2-20221.    Patent Document 3: Japanese Patent Application “Kokai” No."}
{"text": "It is well known in the art to have cigarette holders and ash trays which are incorporated into the dashboard and/or the arm rests of a vehicle. However, with these devices the cigarette smoke is maintained in the vehicle permeating the upholstery and the clothing of the occupants. In addition, when non-smokers are in a vehicle smokers often have to refrain from smoking or hang the cigarette out the window.\nIt would be desirable, therefore, for smokers to have a device for holding a cigarette which directs the smoke out of a vehicle. It would further be desirable, to have a device which may be used on either side of a vehicle. It would further be desirable, to have a device which is removably connectable between a window and door frame of a vehicle."}
{"text": "As the rotor of an electric generator rotates, its magnetic field revolves in unison with the rotor. This rotating magnetic field exerts a cyclical force on the generator's stator core. This force, in turn, causes vibrations to occur in the core which are of concern in high speed, high power machines. The stator core of a machine having a two-pole rotor experiences an elliptical deformation that follows the rotation of the rotor during operation, with resulting vibrations that are at a frequency twice that of the rotational speed of the rotor. The cyclical vibrations will be transmitted through the frame to the generator foundation. The magnitude will depend upon the degree of isolation incorporated in the design.\nIt is known in the art to utilize spring plates to resiliently support a stator core and stator core frame. For example, U.S. Pat. No. 2,320,843, issued to Baudry on Jun. 1, 1943, incorporated by reference herein, discloses a stator support structure made from flat plates mounted tangentially to the perimeter of a stator core frame. The support structure of Baudry includes a pair of vertical spring plates mounted to the perimeter of the stator core frame along opposed sides of the frame. These vertical spring plates are placed into compression by the dead weight of the generator stator and frame. Due to the geometry of the plates, they are rigid in the vertical (compression) direction, but are relatively flexible in the horizontal (bending moment) direction. Additional horizontal spring plates are attached to the perimeter of the stator core frame along its bottom. Due to the geometry of these plates, they provide rigid support in the horizontal (compression) direction, but are relatively flexible in the vertical (bending moment) direction. Together, the vertical and horizontal plates provide rigid vertical and horizontal support for the generator, while at the same time providing a degree of flexibility to accommodate radial vibrations of the stator core. During its operation, an electric generator may experience a short circuit or other abnormal operating condition. Although the prior art support design performs well in isolating the radial vibrations generated during normal operation of the generator, it is less effective in reducing the tangential torque transmitted to the foundation during short circuit conditions. During these irregular operating conditions, an electric generator may experience torque loads as high as twenty times, or more, the normal operating torque. This large torque amplification mandates a very robust support structure, which in turn tends to reduce the desired flexibility in the radial direction.\nThe object of the present invention, therefore, is to provide an electric generator having a support system that is capable of providing the necessary support in both the horizontal and vertical directions while providing sufficient flexibility in the radial direction to dampen rotating force loads, and further to provide sufficient strength and flexibility in the tangential direction to accommodate short circuit and off-normal operating torque loads. It is a further object of this invention to provide an electric generator having a support system that functions to isolate the foundation from the radial vibration forces generated during normal operation and from the tangential forces generated during short circuit and other abnormal operating conditions."}
{"text": "The present invention relates to an apparatus and methods for transferring money or monetary value. More specifically, the present invention relates to a system for managing data to perform send, receive and payment transactions for money transfers.\nIt is increasingly common for funds to be transferred electronically. This can occur from individual to individual or from business to business or between an individual and a business. The transfers may occur within one country or across borders, from one country into another, and thus may involve a currency change.\nIn some cases one or both of the individuals or businesses involved may lack an account at a bank or a credit card. Thus, the funds for the transfer often must be provided in cash (or equivalent) and the payout to the recipient may also preferably be in cash (or equivalent).\nLarge bank and non-bank financial institutions are involved in money transfers. To extend their geographical coverage, the institutions engage agents and may equip them with specialized terminals or other facilities to facilitate communications necessary for money transfers. The personnel of such agents may not be well educated or highly trained and may do several other tasks (such as service as retail sales clerks). This requires a system that is easy to use, requires a minimum of equipment, deters mistakes and fraud and minimizes agent personnel time.\nOne method of initiating or staging such transactions has been the use of a call center that receives calls from persons wishing to transfer money or to receive money transferred. The call center representative collects information sufficient to stage a transaction and provides the caller with a reference or retrieval number that the caller then takes to an agent to fulfill the staged transaction, whether it is a send, receive, bill payment or other transaction. The call center transaction staging interaction results in a staged transaction record that facilitates the fulfillment agent's work to finalize the transaction.\nWhile the call center can provide a personal touch to the service, its staffing can be expensive and may require personnel with special language skills in several languages. In addition, such personnel require additional training when procedures and systems are changed.\nThe agent who is involved in fulfillment may have a compact custom terminal specifically designed for transactions of this type or the agent may have more generalized point of sale device, such as an electronic cash register, equipped with a scanning device for product barcodes, for check reading or optical character recognition or for magnetic card reading, which has associated software for money transfer transactions. Generally, to fulfill a staged transaction, the agent is required to call up and review a record and follow instructions associated with a particular transaction. The agent may enter information to call up the transaction record and to provide new or updated transaction details. This takes time and involves the opportunity for error in manual inputting of information.\nTo make money transfers efficient and to handle the large volume, electronic systems may be used and it is desirable that personnel time is minimized. Thus, there is a need for improved money transfer systems and methods that assist the parties involving in complying with the applicable regulations in various jurisdictions that provide procedures to minimize the amount of personnel time per transaction and that are otherwise economically efficient."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the variety designated CV181138, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include greater yield, better stalks, better roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen comes to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants each heterozygous at a number of loci produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants and have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a boot for downhill skiing.\nIn a generally known manner, such downhill boots have a shell base and an upper.\nThe present invention relates more particularly to a rear entry ski boot, i.e., a boot having a rear portion referred to as a rear spoiler which is journaled on the shell base around a transverse axis. The front portion of the upper can be formed out of a single piece with the shell base or be a separate piece referred to as a cuff which is journaled on the shell base around a transverse axis.\nThe cuff and the rear spoiler are, depending upon the model, generally journaled around the same transverse axis or around two different axes, while for certain boot models, the rear spoiler is journaled on a transverse axis of the cuff.\nThe present invention relates to rear entry ski boots which comprise a rear spoiler journaled on the shell base around a transverse axis, regardless of the more specific arrangements which have been just described.\nThe present invention relates in a more specific manner to those boots which comprise an adjustment apparatus activated by a control element which must be tensioned during closure of the boot before skiing.\nAn adjustment apparatus of the aforementioned type can, for example, be adapted to assure the maintenance of the foot or of the lower leg of the skier or be adapted to limit the flexion of the upper of the boot.\nIt has previously been proposed to link the control element of such an apparatus to the rear spoiler of the boot such that it be tensioned by the movement of the rear spoiler with respect to the shell base during closure of the boot. These systems have the disadvantage of a relatively complex construction and are generally raised in relief with respect to the general profile or shape of the boot, which should be avoided."}
{"text": "Fluoroscopy images, for example, can be acquired with the aid of a C-arc fluoroscopy apparatus. The images then can be transferred to a navigation system, which uses the images as a basis for image-assisted or image-guided surgery. One pre-condition for accurate image-assisted or image-guided surgery is establishing a spatial correlation between a virtual representation of a patient or body data and an actual anatomy of the patient. This process is called registration. The present invention allows the accuracy of this registration process for fluoroscopy images to be verified.\nVerification is beneficial because errors, which are undesirable, can occur in the registration process. Verification detects these errors and, thus, provides an indication of being within a specified or desired tolerance. Registration is often achieved by mapping points or markings on the fluoroscopy images, said markings being arranged in a known position on an attachment of the C-arc radiation source. If, however, an insufficient number of markers are mapped, or the mapping is inaccurate, registration errors can result.\nU.S. Pat. Nos. 6,697,664, 6,675,040 and 6,662,036 include descriptions of the principle of fluoroscopic navigation. However, these references do not provide solutions for reducing or minimizing registration errors."}
{"text": "For laying of offshore pipelines, the commonly known technique is referred to as the S-lay technique. In the S-lay technique, pipe sections are added to the previously launched pipeline in a horizontal position on a lay barge. The pipeline then curves over the stern of the barge, angles down toward the seabed, and curves back to horizontal and lays on the seabed. The profile of the pipeline from the lay barge to the seabed is in the form of a long “S”. Although S-lay technique has been widely utilized, the technique is limited to certain conditions, such as to limited water depth.\nAs an alternative, the J-lay technique has been developed. The J-lay technique is especially advantageous for deep water pipe laying. In the J-lay technique, pipe sections are added to the previously launched pipeline in an essentially vertical position. The pipeline runs down toward the seabed, curves back to horizontal, and lies on the seabed. The profile of the pipeline from the lay barge to the seabed is in the form of a long “J”. The definition of deep water, when referring to the use of J-lay technique, is a direct function of pipe diameter. This relationship is a result of the minimum water depth required for pipe of a certain diameter to achieve the proper flex from the vertical to the horizontal during the vertical laying operation.\nMany different designs for J-lay technique apparatus are known in the current art. Many of the designs include a pivotable tower construction with a hoist system and a pipeline hoist cable suspended from the top of the tower. At the lower end of the tower, a static clamp is provided for holding the previously launched pipeline. The tower is provided with a mobile clamp that can engage the upper end of a new pipe section. The lower end of a new pipe section is connected to the previously launched pipeline, so that during a controlled lowering of the mobile clamp (after the static clamp has been released), the pipeline is lowered and a further new pipe section can be connected to the pipeline. Since the weight of the pipeline in deepwater situations can be enormous (several hundreds of tons), an extremely sturdy design of the tower is required.\nWilkens U.S. Pat. No. 5,464,307 discloses a J-lay technique system, wherein the tower extends downwardly from the deck of the vessel, essentially into the water. The static clamp in the Wilkens reference is provided at the lower end of the tower. The mobile clamp is movable up and down along a track formed in the tower and arranged above the static clamp. In another embodiment of Wilkens a stinger can extends downwardly from the static clamp. The static clamp in the Wilkens reference is mounted at deck level. The stinger extends into the water. The mobile clamp engages on the pipeline at a position below the static clamp and is guided along the stinger.\nIn many drilling techniques, vessels, particularly semi-submersible drilling vessels, are not used efficiently since drilling activities are not always needed while the vessel is under contract. Also other types of vessels are not used efficiently because those are designed for only one purpose. The Wilkens reference discloses an attempt to modify drilling vessels into pipe laying vessels. This was not successful.\nOther references incorporated into this application include de Varax U.S. Pat. No. 6,361,250, and Seguin U.S. Pat. No. 6,352,388.\nA need exists for an object to provide an improved pipelaying system for vertical or near vertical marine pipe laying, in particular in deepwater. A need exists for a (semi-submersible) drilling vessel that can be used effectively as pipe laying vessel for the vertical or near vertical laying of offshore pipelines. A need also exists for a floating production vessel that can be used effectively also as pipe laying vessel for the vertical or near vertical laying of offshore pipelines.\nA need exists for a more efficient method for the vertical or near vertical laying of offshore pipelines."}
{"text": "The present invention relates to bonding pads for semiconductor chips.\nSemiconductor chips have bonding pads which are relatively large areas of metalization providing a contact area for an assembler to solder a wire connection to one of the pins of an integrated circuit package. This wide metal area of the bonding pad is then coupled to the tiny metal connections of the integrated circuit. The semicustom method of designing and fabricating integrated circuits uses prefabricated semiconductor wafers which have completed circuits except for the final process step of device or component interconnection. A cost savings is realized by adapting a standard wafer, which has already been designed and manufactured, to a particular customer's needs by simply designing appropriate metal interconnections. Thus, for each customer, a metal mask is designed which connects the necessary active and passive semiconductor components to form the circuit configuration required by the customer. This method of design and fabrication allows development of a semicustom circuit in less time and at lower cost than if the entire wafer had to be independently designed for each customer.\nOne limitation in semicustom design is that the number of bonding pads required to form connections to a particular customer's circuit will vary from customer to customer. Accordingly, a standard chip design will have enough bonding pads to accommodate the largest number anticipated to be required by any customer. This results in a number of bonding pads being wasted space for a large number of customers. A bonding pad will typically be a cell area which is simply an N-doped epitaxial layer on top of a P-doped substrate. Inherent in this structure is a clamping diode or a capacitance. For the clamping diode, the N portion of the diode is the epitaxial region and the P portion of the diode is the substrate. The junction capacitor is formed at the junction of the substrate and the epitaxial region. Thus, in some circuit designs, it may be possible to take advantage of this inherent capacitance or clamping diode."}
{"text": "Pressure gauges are widely utilized in a myriad of different environments for commercial and industrial applications. Typically, pressure gauges measure pressure and provide an indication of the value. Pressure values are usually displayed in analog form (e.g., by a pointer) or in digital form (e.g., by an LED readout). Gauges displaying values in analog form often include a mechanical pressure transducer, such as a Bourdon tube, that displaces a predictable amount in response to the pressure to which the transducer is exposed. The displacement is translated via a movement to a rotatable pointer that moves opposite a dial of calibrated pressure values. Gauges displaying values in digital form often use an electronic pressure transducer, such as a piezosensor, that generates an electrical signal in response to the pressure to which the transducer is exposed. The electrical signal is then converted to a symbol that is presented on a display."}
{"text": "1. Field of the Invention\nThe present invention relates to slime control agents especially useful in the paper and pulp industries, more particularly it relates to a slime inhibitor containing as the effective component at least one .alpha.-chloro-o-acylbenzaldoxime derivative represented by the following general formula ##SPC1##\nWherein Y represents a halogen atom, a nitro group, a lower alkyl group, a lower alkoxyl group, a lower alkylcarboxy group, or a lower alkylamino group; R represents lower alkyl group or a halogenated alkyl group; and n represents O or an integer of 1-3.\nThe invention further relates to a slime control composition containing the aforesaid .alpha.-chloro-o-acylbenzaldoxime derivative and a known slime control agent such as a halogenated phenol or methylene bisthiocyanate as effective components, and also relates to a method of controlling slime using all such compositions.\n2. Description of the Prior Art\nIn paper manufacturing, microorganisms are apt to generate quite easily in various processing streams. In such a case it sometimes happens that the microorganisms and the secretions thereof form a viscous material called slime on chests, pipes, and the like during the manufacture of paper. The generation of such slime is accompanied by a reduction in working efficiency, tearing of paper and the formation of stains in products, and thus causes a great amount of trouble in the production of paper and pulp. Therefore, the occurence of slime must be prevented in the manufacture of paper and pulp and it is necessary to kill the slime-forming microorganisms and to restrain the generation of such microorganisms during the manufacture of paper or pulp.\nAs is well known to those skilled in the art, the slime produced in the paper mill is a viscous mass having various forms and appearances such as a soft paste-like, gelatinous, gummy, fibrous, hard or horny mass which emits unpleasant odors. In many cases, the slime is a colored amorphous mass and is generally composed of heterogeneous materials, i.e., microorganisms, such as bacteria, fungi, yeast and algae and the metabolites thereof, pulp fibers and various additives incorporated into the pulp stream. The problems caused by slime are usually due to a mass as described above and are not serious if the heterogeneous materials are present in the discrete state even though the materials are contained in a relatively large amount in the pulp. A mass of slime is considered to be formed by growing viscous masses as nuclei produced by pulp fibers, additives, microorganisms and the like where the inner wall of paper mill system is rough and where the flow stream of the pulp or pulp suspension stagnates. The microorganisms, main source of the slime, are usually borne by industrial water consumed in a large quantity and raw pulp woods into paper processing systems. Additionally, since most paper processing steps are conducted while being exposed to air, there is the possibility that the pulp will be contaminated with various microorganisms floating in the air. Once the pulp is contaminated with the microorganisms, they usually grow abundantly in the raw pulp since the raw pulp and the white liquor can offer optimum conditions for microorganism growth in view of their nutrients, moisture, temperature, pH and oxygen supply. The slime formed in the paper mill systems causes various colored spots or stains in the paper product and may be the main factor in a decrease in strength and an increase in breaking of the paper product.\nThe slime also causes contamination and/or clogging of wires and felts used in paper processing and adversely affects the durability of such equipment due to deterioration of the material of the equipment. Thus, the operation of pulp and paper processings must be discontinued for cleaning, thereby decreasing the production efficiency and causing damage to the raw pulp materials.\nIt is advantageous to control or kill microorganisms having a tendency to produce such slime by a chemical treatment, i.e., by incorporating various slime controlling agents into the process streams. Slime is generally produced by various microorganisms such as bacteria, fungi, yeast and the like, so it is usually necessary to use different types of slime controlling agents at the same time. In light of the above, it would be a great advantage in the art if various microorganisms grown in the process streams in a paper mill could be controlled by a single chemical agent.\nThe slime controlling agents which have conventionally been used are organometallic compounds, chlorinated phenols and dithiocarbamate derivatives. These well known compounds are not satisfactory because they are toxic to humans, emit an unpleasant odor per se or even in the paper products, lower the whiteness of the paper, decrease the sizing efficiency, etc."}
{"text": "The present invention relates to upright vacuum cleaners and, more particularly, to an upright vacuum cleaner having a floor cleaning nozzle and an above-floor cleaning nozzle and two motor and fan units as vacuum sources.\nUpright vacuum cleaners that utilize two motor and fan units for providing sources of suction are known where one motor and fan unit draws air through the walls of a dirt collection bag, or the like, commonly called a clean air motor, and a second motor and fan unit is disposed in the cleaning head of the vacuum cleaner through which dirt laden air passes, commonly called a dirty air motor. However, some of these prior art vacuum cleaners turn one motor and fan unit off when the other is on since they are used for different purposes. Typically, the clean air motor and fan unit is used for off the floor cleaning with accessories for cleaning furniture and draperies and the like, whereas the dirty air motor and fan unit is used for floor cleaning, such as disclosed in U.S. Pat. No. 4,225,999. In addition, in some known prior art vacuum cleaners where there is utilized simultaneously both motor and fan units, there is also utilized a third motor for driving the brush roller in the cleaning head, such as is disclosed in U.S. Pat. No. 5,134,752. In known prior art vacuum cleaners where there is utilized simultaneously both motor and fan units the operator cannot selectively turn one of the motor and fan units off, if desired, for cleaning different surfaces."}
{"text": "Musical performances have become a key component of electronic and multimedia products such as stand-alone video game devices, computer-based video games, computer-based slide show presentations, computer animation, and other similar products and applications. As a result, music generating devices and music playback devices are now tightly integrated into electronic and multimedia components.\nMusical accompaniment for multimedia products can be provided in the form of digitized audio streams. While this format allows recording and accurate reproduction of non-synthesized sounds, it consumes a substantial amount of memory. As a result, the variety of music that can be provided using this approach is limited. Another disadvantage of this approach is that the stored music cannot be easily varied. For example, it is generally not possible to change a particular musical part, such as a bass part, without re-recording the entire musical stream.\nBecause of these disadvantages, it has become quite common to generate music based on a variety of data other than pre-recorded digital streams. For example, a particular musical piece might be represented as a sequence of discrete notes and other events corresponding generally to actions that might be performed by a keyboardist—such as pressing or releasing a key, pressing or releasing a sustain pedal, activating a pitch bend wheel, changing a volume level, changing a preset, etc. An event such as a note event is represented by some type of data structure that includes information about the note such as pitch, duration, volume, and timing. Music events such as these are typically stored in a sequence that roughly corresponds to the order in which the events occur. Rendering software retrieves each music event and examines it for relevant information such as timing information and information relating the particular device or “instrument” to which the music event applies. The rendering software then sends the music event to the appropriate device at the proper time, where it is rendered. The MIDI (Musical Instrument Digital Interface) standard is an example of a music generation standard or technique of this type, which represents a musical performance as a series of events.\nComputing devices, such as many modern computer systems, allow MIDI data to be manipulated and/or rendered. These computing devices are frequently built based on an architecture employing multiple privilege levels, often referred to as user-mode and kernel-mode. Manipulation of the MIDI data is typically performed by one or more applications executing in user-mode, while the input of data from and output of data to hardware is typically managed by an operating system or a driver executing in kernel-mode.\nSuch a setup requires the MIDI data to be received by the driver or operating system executing in kernel-mode, transferred to the application executing in user-mode, manipulated by the application as needed in user-mode, and then transferred back to the operating system or driver executing in kernel-mode for rendering. Data transfers between kernel-mode and user-mode, however, can take a considerable and unpredictable amount of time. Lengthy delays can result in unacceptable latency, particularly for real-time audio playback, while unpredictability can result in an unacceptable amount of jitter in the audio data, resulting in unacceptable rendering of the audio data.\nThe invention described below addresses these disadvantages, providing kernel-mode audio processing modules."}
{"text": "In general, today's advanced material applications are subjected to environments and stresses, which benefit from combinations of material properties. For example, in ballistic applications, a material is sought which is lightweight and thus fuel efficient, while at the same time provides great impact absorption properties to prevent injury or mechanical failure to an underlying structure that may be the target of shrapnel or an exploding device. In aircraft or seacraft applications, materials that are strong, light-weight and at the same time corrosion resistant are also sought. In an attempt to achieve these and other material property combinations, composite materials (i.e., multiphase materials) are employed.\nThere are many types of composite materials. For example, particle-reinforced composite materials, fiber-reinforced composite materials, structural composite materials or layered composite materials are generally well-known. Each type of composite material can include two or more phases wherein one phase makes up the majority of the material and is know as the matrix material and the second phase (and potentially additional phases) make(s) up a lesser extent of the composite and can be dispersed within the matrix material or layered within the matrix material to form a sandwich. The presence of the second and additional phases affects the material properties (such as, for example, the mechanical and thermal properties) of the composite material. That is, the material properties of the composite material are dependent upon the material properties of the first phase and the second phase (and additional phases) as well as the amounts of the included phases forming the composite. Thus, material properties of a composite can be tailored for a specific application by the selection of specific concentrations of the phases, as well as potentially, the sizes, shapes, distribution, and orientation of the included phases.\nDifficulties in the formation, durability, and tailoring of material properties have however impeded or prevented the use of composite materials in some applications. For example, material failure may be due, at least in part, to abrupt property changes along phase interfaces."}
{"text": "This is a continuation-in-part of co-pending patent application Ser. No. 07/198,384 filed 1/25/88, now U.S. Pat. No. 4,883,103, issued Nov. 28, 1989, which is a continuation-in-part of application Ser. No. 106,632, filed 10/8/87, now U.S. Pat. No. 4,816,045, issued Mar. 28, 1989, which is a continuation-in-part of application Ser. No. 896,081 filed 3/31/86, now U.S. Pat. No. 4,770,677, issued Sept. 13, 1988.\nThe present invention relates to a sealing assembly in a filler neck, and particularly to means for establishing a seal between the filler neck and a liquid-dispensing nozzle in the filler neck. More particularly, the present invention relates to a sealing assembly in a vehicle fuel system filler neck for sealingly admitting a fuel-dispensing nozzle into the filler neck without permitting substantial exhaustion of fuel vapors in the filler neck to the atmosphere or admission of atmospheric pressure into the fuel system through the filler neck.\nPolluting emissions are released into the atmosphere during each and every motor vehicle refueling activity through the mouth of the filler neck. In view of the frequency of vehicle refueling activities, it will be appreciated that a significant quantity of otherwise combustible fuel is lost to the atmosphere each time a vehicle is refueled.\nNew fuel vapors are generated during refueling due to splash and agitation of the dispensed fuel, as well as from potential temperature differences between the fuel tank and the dispensed fuel. In addition, as the liquid fuel dispensed at the pump fills the vehicle fuel tank, fuel vapors that are present in the tank and generated during refueling are displaced by liquid fuel. These displaced fuel vapors are moved out of the fuel tank vapor space by the displacing action of the liquid fuel. In conventional vehicle fuel systems, these displaced vapors are released directly into the atmosphere via the fuel tank filler neck and are a contributing factor to air pollution.\nMany conventional filler necks and fuel tank assemblies on vehicles are not configured to control this escape of fuel vapor effectively during the refilling process. Generally, the filler neck of a vehicle has a diameter greater than the diameter of the fuel-dispensing nozzle. Because of the relatively large diameter of the filler neck with respect to the fuel-dispensing nozzle, the fuel vapor created by fuel dispensed into the fuel tank through the filler neck escapes out through the filler neck and to the atmosphere substantially continuously during the filling operation. This creation of fuel vapor during the refueling process has recently been recognized as a major contributor to pollution of the atmosphere, especially with newer, more volatile blends of vehicle fuel.\nAttempts have been made to control this escape of fuel vapor during the filling process by adapting the service station fuel-dispensing nozzle to include a vapor recovery portion which captures the escaping fuel vapor and prevents the fuel vapor from escaping into the atmosphere. Such conventional methods of capturing the escaping fuel vapor have generally proven unsatisfactory because of the complexity of the system and because of relatively poor vapor-capturing characteristics of the system. It is generally recognized that for satisfactory fuel vapor control during refilling, a system must be provided which is on-board the vehicle for capturing and controlling such fuel vapor generated during refilling. Ideally, such a system would work automatically without any operator supervision.\nAnother problem with many conventional filler neck and fuel tank assemblies is that it is sometimes extremely difficult to achieve accurate filling levels because of poorly designed filler neck seals which leak, resulting in an inability to maintain any vacuum that might develop in the filler neck itself during introduction of fuel into the filler neck. Conventional fuel nozzles are often equipped with vacuum-actuated systems that shut off flow of fuel through the nozzle upon exposure of a sensor in the nozzle to a predetermined negative pressure. Typically, the splashed back fuel enters and blocks a conduit formed in the nozzle to lie in communication with the sensor to generate the negative pressure, thereby shutting off fuel flow.\nRelying upon fuel splash-back provides generally unsatisfactory fill level consistency because such fuel splash back is dependent upon a number of factors. Such factors include the flow rate of the nozzle, the configuration and routing of the filler neck, and the general orientation of the filler neck with the fuel tank, as well as fuel-dispensing nozzle sensitivity. Because of such inconsistency in achieving accurate fill levels, it would be advantageous to provide a nozzle shutoff control system having an improved means for providing a negative pressure signal to a vacuum-actuated shutoff sensor which would ideally permit the fuel tank to be filled to accurate, consistent fill levels during each refueling operation.\nOne object of the present invention is to provide a sealing assembly for partitioning a fuel-conducting passageway in a filler neck to limit flow of vapors and gases through the passageway during dispensing of fuel into the filler neck and past the sealing assembly.\nAnother object of the present invention is to provide a sealing assembly for preventing fuel vapor loss to the atmosphere through the filler neck during refueling in cooperation with a fuel vapor recovery system.\nYet another object of the present invention is to provide a sealing assembly for blocking communication of an inner chamber of the filler neck with the atmosphere during refueling to maintain a negative pressure condition in the inner chamber in cooperation with a vapor-handling system to enhance operation of a vacuum-actuated nozzle shutoff system communicating with the inner chamber.\nAnother object of the present invention is to provide a filler neck sealing assembly having an elastic seal configured to maintain unbroken sealing engagement with the nozzle during insertion of the nozzle into and withdrawal of the nozzle from the filler neck, even though the nozzle may have been deformed to assume an out-of-round shape.\nYet another object of the present invention is to provide a filler neck sealing assembly having an elastic seal that is configured to be stretched or distended easily to conform to the contours of an out-of-round fuel-dispensing nozzle.\nAccording to the present invention, a filler neck seal assembly is provided for use in a filler neck sized to receive a fuel-dispensing nozzle. The seal assembly includes an elastic seal member and a seal retainer. The seal member has an annular fuel nozzle-sealing portion and an annular filler neck-mounting portion which surrounds the annular fuel nozzle-sealing portion. The seal retainer includes means for engaging a filler neck containing the seal member and means for unyieldingly clamping the annular filler neck mounting portion of the seal member to the filler neck to establish a fixed circumferentially extending liquid fuel and fuel vapor seal between the seal member and the filler neck. The filler neck-mounting portion of the seal member is clamped to cantilever the annular fuel nozzle-sealing portion in the filler neck for movement relative to the clamped filler neck-mounting portion in response to engagement with a fuel nozzle during movement of the fuel nozzle in the filler neck.\nIn preferred embodiments, the annular fuel nozzle-sealing portion includes an annular, radially inwardly facing, perimeter end wall defining a nozzle-receiving aperture and inclined annular outer and inner sealing faces. The outer sealing face is presented in an axially outwardly direction toward the filler neck mouth to engage the nozzle in sealing relation during insertion of the nozzle into the nozzle-receiving aperture. The inner sealing face is presented in an axially inwardly facing direction away from the filler neck mouth to engage the nozzle in sealing relation during withdrawal of the nozzle from the nozzle-receiving aperture. The outer and inner sealing faces are interconnected by the annular, radially inwardly facing, perimeter end wall and diverge in a radially outwardly extending direction away from the perimeter end wall and toward the filler neck-mounting portion.\nThe elastic annular seal member is formed to include an axially outwardly facing first annular groove and an axially extending second annular groove between the fuel nozzle-sealing and filler neck-mounting portions. A radially extending, axially outwardly facing surface defines a bottom wall of the first annular groove and a radially extending, axially inwardly facing surface defines a bottom wall of the second annular groove. The filler neck includes positioning means extending into the second annular groove formed in the seal member for positioning the filler neck-mounting portion in sealing engagement with the filler neck.\nThe seal retainer includes retaining means extending into the first annular groove formed in the seal member for retaining the filler neck-mounting portion in sealing engagement with the filler neck to establish fixed circumferentially extending liquid fuel and fuel vapor seal between the seal member and the filler neck. The seal retainer includes means for rotatably engaging the filler neck to place the retaining means in the first annular groove formed in the seal member.\nOne feature of the present invention is the provision in a filler neck of seal-establishing means for admitting a fuel-dispensing nozzle into an inner chamber in the filler neck without coupling the inner chamber to the atmosphere through the filler neck mouth. One advantage of such a sealing system is that it helps to prevent fuel vapor loss through the filler neck during operation of an on-board vehicle fuel vapor recovery system by blocking flow of pressurized fuel vapor in the filler neck to the atmosphere during refueling. Another advantage of this sealing system is that it cooperates with a vapor-handling system in a vehicle fuel system to maintain any negative pressure condition that develops in the inner chamber during refueling to enhance the operation of a vacuum-actuated nozzle shutoff system having its shutoff sensor communicating with the inner chamber.\nAnother feature of the present invention is that the seal member includes a radially outer portion that is securely clamped in place in the filler neck to establish a fixed circumferentially extending liquid fuel and fuel vapor seal in the filler neck between the seal member and the filler neck. Advantageously, such a fixed seal is less likely to be broken as a result of wear and tear on the seal member caused by repeated insertion and removal of the fuel-dispensing nozzle during refueling.\nStill another feature of the present invention is that the seal member includes a radially inner portion that is configured to stretch and move relative to the clamped radially outer portion when engaged by the fuel-dispensing nozzle so that the seal member is stretched or distended to conform to the contour of a conventional or out-of-round nozzle during insertion of the nozzle into the filler neck and withdrawal of the nozzle from the filler neck. In a preferred embodiment, the seal member includes a thin annular web that is configured to interconnect the relatively thicker unclamped nozzle-sealing portion and clamped seal-mounting portions of the seal member. Provision of a thinner annular web allows the movable nozzle-sealing portion to stretch more easily to follow a regular or irregular contour of the fuel-dispensing nozzle as it is moved within the filler neck during refueling.\nAdditional objects, features, and advantages of the invention will become apparent to those skilled in the art upon consideration of the following detailed description of preferred embodiments exemplifying the best mode of carrying out the invention as presently perceived."}
{"text": "1. Field of the Art\nThe invention relates to apparatuses for the generation of power which comprise at least a Diesel engine having a plurality of sequentially operating cylinders which are supercharged by a plurality of alternating air compressor cylinders which deliver scavenging air to the engine.\n2. Prior Art\nMost existing Diesel engines which are supercharged by an alternating compressor having a plurality of cylinders which operate sequentially include a manifold which feeds all engine cylinders. All compressor cylinders open into that collector. As a consequence, the engine cylinders are scavenged under an inlet pressure which is approximately constant. A thermo-dynamic study of scavenging indicates that the work consumed for compression and air transfer is directly related to the difference between the mean inlet pressure of the scavenging air and the outlet pressure. As a consequence, the transfer work is determined by the mean pressure in the collector rather than by the outlet pressure of a particular compressor cylinder. However, if the apparatus has a single motor cylinder fed by a single compressor cylinder (as described for instance in French Pat. No. 1,238,426), the situation is completely different since there is no manifold which feeds a plurality of motor cylinders."}
{"text": "As the frequency and circuit density of an electronic device become higher and its size becomes smaller, a problem related to interference due to radiation of an undesired electromagnetic wave arises. It is therefore required to suppress leakage of an undesired electromagnetic wave to the outside. There is known a method of covering, with a conductive resin film, the surface of a nonconductive resin layer for sealing a semiconductor chip in order to impart a shielding function to a semiconductor package. There is also proposed a technique of realizing a module which includes a transmission/reception antenna and has a shielding function by forming apertures in a portion, of a conductive resin film and a nonconductive resin layer for sealing a semiconductor chip, which covers the upper surface of the semiconductor chip.\nIn the conventional techniques, the surface of the semiconductor package with the antenna is covered with a conductive resin film, and apertures for enabling radiation and reception of a desired electromagnetic wave to be used for communication are formed in the conductive resin film. In this case, a diffraction wave occurring in an edge portion of the conductive resin film degrades an antenna gain in a desired radiation direction depending on the package size and the frequency of the desired electromagnetic wave."}
{"text": "a) Field of the Invention\nThis invention relates to unsymmetrical straight chain alkyl-substituted phenylxanthines which demonstrate potent bronchodilating activity. The compounds are 8-phenylxanthines having a methyl or propyl group at the 1- and 3- positions with the proviso that the 1- and 3-substituents may not be the same and the pharmaceutically acceptable salts of such compounds.\nb) State of the Art\nXanthines of various types have been used or proposed as drugs for various indications. For example, theophylline and aminophylline relax the smooth muscle of the bronchial airways and pulmonary blood vessels, thereby acting as pulmonary vasodilators, bronchodilators and smooth muscle relaxants. Like other xanthines these compounds possess the following actions as well: coronary vasodilator, diuretic, cardiac and cerebral stimulant and skeletal muscle stimulant. Dyphylline is another xanthine having activity similar to that of theophylline and aminophylline. Most of these xanthines bear identical alkyl substituents in positions 1 and 3. Nevertheless, a large number of xanthines having different substituents in these positions are known [see, e.g., R.F. Bruns, G.H. Lu, and T. Pugsley, Mol. Pharmacol., 29, 331 (1986); J.W. Daly, W.L. Padgett, and M.T. Shamin, J. Med. Chem , 29, 1305 (1986); J.N. Wells, J.E. Garst, and G.L. Kramer, J. Med. Chem., 24, 954 (1981)] and several have notable pharmacological activity. For example, enprofylline (3-propylxanthine) has bronchodilating and antiasthmatic characteristics [see, e.g., C.G.A. Persson, K.-E.-Andersson and G. Kjellin, Life Sci., 38, 1057 (1986)]and IBMX (3-isobutyl-1-methylxanthine) and related compounds are potent inhibitors of phosphodiesterase [M. Chasin and D.N. Harris, Adv. Cyclic Nucleotide Res., 7, 225 (1976)].\nNevertheless, relatively little prior art exists for 8-phenylxanthines with an unsymmetrical alkyl substitution of the 1 and 3 positions. Bruns et al. [Proc. Natl. Acad. Sci. USA, 80, 2077 (1983)]disclose many 8-arylxanthines. One preferred compound is 1,3-dipropyl-8-(2-amino-4-chlorophenyl)xanthine (PACPX). The only unsymmetrical alkyl-substituted xanthines which Bruns et al. disclose lack an 8-phenyl substituent.\nMany 8-arylxanthines are disclosed in U.S. Pat. No. 4,593,095 to Snyder et al., including the preferred compound PACPX. As with Bruns et al., most of the Snyder et al. compounds are symmetrically substituted. Although one of the preferred compounds of the present invention is generically within Snyder's formula, none is exemplified nor claimed. Indeed, the only exemplified unsymmetrically substituted 1,3-dialkyl-8-phenylxanthine is 1-allyl-3-methyl-8-phenylxanthine [U.S. Pat. No. 4,593,095, column 11; see, also, H.W. Hamilton et al., J. Med. Chem., 28, 1071 (1985)].\nOther unsymmetrically substituted 1,3-dialkyl-8-phenylxanthines described in the literature are 1-isoamyl-3-isobutyl-8-phenylxanthine [J.W. Daly et al., J. Med. Chem., 28, 487 (1985)], 3-ethyl-1-methyl-8-phenylxanthine [H.W. Hamilton et al., J. Med. Chem., 30, 91-96 (1987)], and 1-hexyl-3-methyl-8-phenylxanthine [H.G.V. Schub, German Patent 1,091,570 (1960)]. The first two compounds have been studied only in binding assays for affinity for the A.sub.1 and A.sub.2 subtypes of adenosine receptors. For the third compound \"drug use\" is claimed. The possibility that the therapeutic responses to adenosine receptor antagonists, such as theophylline, may involve such antagonism is uncertain. This uncertainty is based largely on the observation that enprofylline, which is about five-fold more potent than theophylline in man and other species [C.G.A. Perrson, Agents Actions, 13 (Suppl.), 115-129 (1983)], is much less effective than theophylline in diminishing responses to adenosine in all tissues studied, except for the rat hippocampus [B.B. Fridholm and C.G.A. Persson, Eur. J. Pharmacol., 81 673-676 (1982)]. Thus, it is not certain to what extent the two xanthines share common cellular actions or whether adenosine receptor blockade is of significance in the bronchodilator effect that theophylline produces in asthmatic patients [T.W. Rall in \"Goodman and Gilman's, The Pharmacological Basis of Therapeutics,\" 7th Ed., A.G. Gilman, L.S. Goodman, T.W. Rall and F. Murad, eds., MacMillan Publishers, New York 1985, pp. 589-603]. Thus, the potent bronchodilating activity of the compounds of this invention is unexpected and the mechanism by which they produce this pharmacological response remains to be elucidated.\nCo-pending application Serial No. 108,990 filed Oct. 1, 1987 (now U.S. Pat. No. 4,783,530) describes various 1,3-alkyl substituted-8-(3,4-or 4-substituted phenyl)xanthines. The 1,3-substituents need not be the same. Among the preferred compounds are: 1-propyl-3-methyl-8-(3-N,N-dimethylaminomethyl-4-hydroxyphenyl)-xanthine and 1-propyl-methyl-8-(4-cyanophenyl-)xanthine.\nThis invention provides a series of novel 1,3-unsymmetrical straight chain alkyl-substituted 8-phenylxanthines which produce potent bronchodilating activity in a test for reversal of histamine-induced bronchoconstriction, the preferred test for measuring this activity (T.W. Rall, loc. cit.), as well as in a test for prevention of antigen-induced bronchoconstriction in guinea pigs. The compounds also have a decreased tendency to produce unwanted side effects, such as increasing heart rate, decreasing blood pressure and causing emesis, as is seen with prototypic xanthine-type bronchodilators, e.g., theophylline. Binding studies demonstrate that this beneficial bronchodilating action is not correlated with A.sub.1 - or A.sub.2 -adenosine receptor binding."}
{"text": "The present invention relates to a drive circuit for driving a display panel such as a liquid crystal display panel, and a display device including the display panel.\nIn a conventional liquid crystal display device of an active matrix type, a liquid crystal display panel generally has a plurality of display pixels arranged in a matrix pattern. Each of the display pixels includes a liquid crystal layer sandwiched between two glass plates and an active element such as TFT (Thin Film Transistor) for controlling a voltage applied to the liquid crystal layer.\nFurther, the conventional liquid crystal display device includes a drive circuit for driving the liquid crystal display panel. The drive circuit includes a gate driver and a source driver. The gate driver is provided for supplying a control signal through a scanning line (a gate line) for tuning on or off each of the active elements. The source driver is provided for supplying a gradation voltage through a data line (a source line) to opposite electrodes sandwiching the liquid crystal layer of each of the display pixels.\nWhen an electrical field is applied to the liquid crystal layers with a direct current with the same polarity for a prolonged period of time, the liquid crystal layers tend to deteriorate. In order to prevent the deterioration, a specific drive method has been widely adopted. In the drive method, the graduation voltage with a positive polarity and the graduation voltage with a negative polarity are alternately supplied to each of the display pixels. The graduation voltage with the positive polarity is greater than a reference voltage, and the graduation voltage with the negative polarity is smaller than the reference voltage.\nFor example, when a dot inversion method is adopted as the drive method, the graduation voltage with the positive polarity and the graduation voltage with the negative polarity are alternately supplied per dot (for example, per display pixel). When a line inversion method is adopted as the drive method, the graduation voltage with the positive polarity and the graduation voltage with the negative polarity are alternately supplied per line.\nWhen the liquid crystal display device adopts the drive method described above, the source driver generally includes an impedance conversion circuit. The impedance conversion circuit includes an operation amplifier (referred to as a high voltage side operation amplifier) for outputting an analog voltage with the positive polarity and an operation amplifier (referred to as a low voltage side operation amplifier) for outputting an analog voltage with the negative polarity.\nPatent Reference 1 has disclosed a conventional liquid crystal display device of an active matrix type. According to Patent Reference 1, a source driver includes a high voltage side operation amplifier (or a positive amplifier) of a voltage follower type and a low voltage side operation amplifier (or a negative amplifier) of the voltage follower type.\nIn the source driver disclosed in Patent Reference 1, the high voltage side operation amplifier is operated with a group of a power source voltage AVDD and a power source voltage AGNDP (AVDD>AGNDP), and the low voltage side operation amplifier is operated with a group of a power source voltage AVDDN and a power source voltage AGND (AVDDN>AGND). Accordingly, four power source voltages AVDD, AGNDP, AVDDN, and AGND are used for driving the high voltage side operation amplifier and the low voltage side operation amplifier.    Patent Reference 1: Japanese Patent Publication No. 2006-292807\nPatent Reference 2 has disclosed another conventional liquid crystal display device. According to Patent Reference 2, a source driver includes a high voltage side operation amplifier operated with a group of a power source voltage VLCD and a power source voltage VMM (VMM=1/2 VLCD) and a low voltage side operation amplifier operated with a group of the power source voltage VMM and a ground GND (GND=0 V). Accordingly, three power source voltages VCLD, AMM, and GND are used for driving the high voltage side operation amplifier and the low voltage side operation amplifier.    Patent Reference 2: Japanese Patent Publication No. 10-062744\nAccording to Patent Reference 1 and Patent Reference 2, the source driver further includes a switch circuit for connecting each of the data lines alternately to an output terminal of the high voltage side operation amplifier and an output terminal of the low voltage side operation amplifier. Accordingly, the analog voltage with the positive polarity and the analog voltage with the negative polarity are alternately supplied to the data lines.\nMore specifically, the switch circuit is configured to connect the data line at an i-th position to the high voltage side operation amplifier and the data line at an (i+1)-th position adjacent to the i-th position to the low voltage side operation amplifier. Accordingly, the analog voltage with the positive polarity is supplied to the data line at the i-th position, and the analog voltage with the negative polarity is supplied to the data line at the (i+1)-th position. At this moment, a potential of the data line at the i-th position is lower than a reference voltage, and a potential of the data line at the (i+1)-th position is higher than the reference voltage.\nAt the next timing, the switch circuit switches the data line at the i-th position from the high voltage side operation amplifier to the low voltage side operation amplifier, and the data line at the (i+1)-th position from the low voltage side operation amplifier to the high voltage side operation amplifier. As a result, the analog voltage with the negative polarity is supplied to the data line at the i-th position, and the analog voltage with the positive polarity is supplied to the data line at the (i+1)-th position. Accordingly, the potential of the data line at the i-th position transits from a high potential to a low potential, and the potential of the data line at the (i+1)-th position transits from a low potential to a high potential.\nIn the conventional liquid crystal display device, when the switch circuit switches the data line, the output terminal of the high voltage side operation amplifier is switched from the data line at the i-th to the data line at the (i+1)-th, and the output terminal of the low voltage side operation amplifier is switched from the data line at the (i+1)-th to the data line at the i-th. Accordingly, a potential of the output terminal of the high voltage side operation amplifier may abruptly decrease, and a potential of the output terminal of the low voltage side operation amplifier may abruptly increase temporarily.\nWhen the potential of the output terminal of the high voltage side operation amplifier abruptly decreases, or the potential of the output terminal of the low voltage side operation amplifier abruptly increases, a bias in a forward direction is applied to a parasite bipolar transistor inside the high voltage side operation amplifier or a parasite bipolar transistor inside the low voltage side operation amplifier. As a result, an excess current may flow in the parasite bipolar transistor, thereby damaging the high voltage side operation amplifier or the low voltage side operation amplifier.\nIn view of the problems described above, an object of the present invention is to provide a drive circuit for driving a display panel, and a display device including the display panel capable of solving the problems of the conventional display device. In the present invention, it is possible to prevent an excess current from flowing in a high voltage side operation amplifier and a low voltage side operation amplifier.\nFurther objects and advantages of the invention will be apparent from the following description of the invention."}
{"text": "One of the problems which is encountered with the conventionally used railway car trucks is the rock and roll encountered due at least in part to the staggered rail joints generally used in the United States.\nAnother problem is that of hunting or the tendency of the truck to oscillate resonantly about the center plate of the car body, due to dynamic instability.\nIt is preferred to solve the problems of rock and roll hunting with a construction wherein the static height of the truck between empty and full loads is not too great and preferably not greater than about 21/2inch. It further is desired that the railway car truck be capable of utilizing a standard wheel base, a standard center plate, and standard brakes.\nIt is an object of the present invention to reduce or eliminate rock and roll commonly occurring in conventional railway car trucks.\nAnother object of the present invention is to reduce or eliminate hunting commonly occurring in conventional railway car trucks.\nAnother object of the present invention is to increase the speed at which the instability that results in hunting occurs to a speed beyond the usual maximum operating speed of freight trains.\nAnother object of the present invention is to provide a truck which reduces wheel flange wear.\nAnother object of the present invention is to provide a truck having a standard wheel base.\nAnother object of the present invention is to provide a railway car truck in which the car may utilize standard brakes.\nAnother object of the present invention is to provide a railway car truck which can utilize a standard center plate.\nAnother object of the present invention is to reduce or eliminate the problem of car center plate failures which is presently commonly occurring.\nAnother object of the present invention is to provide a railway car truck wherein the difference in the static height of the truck between empty and fully loaded car is not too great.\nAnother object of the present invention is to provide a truck design which will allow the attachment of an interface between the car body and truck of a sufficient area to transmit longitudinal loads between the car body and truck without inducing substantial torsional deflections into the transverse member.\nOther objects will be apparent from the following description and drawings."}
{"text": "In recent years the telecommunications industry has increased its demand for improved performance from current controlled oscillators (CCO). For example, when designing phase locked loops (PLL) for frequency synthesizers and clock recovery circuits, it helps to have a CCO with linear gain to allow better modeling during system design. Better modeling during system design helps avoid possible instability problems.\nAdditionally, it is important to reduce the CCO's power consumption and reduce the design margin. This can be achieved by designing the CCO to have low process and temperature sensitivity.\nConventional 3-stage ring oscillators of the prior art can have a wide tuning range, but the CCO gain is sensitive to process and temperature variation. The CCO gain is much higher when it works under low temperature, fast-fast (FFL) conditions than it works under high temperature, slow-slow (SSH) conditions. In order to make a conventional CCO oscillate over a certain frequency range, a much larger tuning range is required because of the process and temperature variations. Another problem with conventional CCO's is that the gain will drop, or become flat, at high frequencies, rather than increasing linearly, because of velocity saturation.\nFIG. 1 illustrates a prior art circuit 7 comprising a conventional CCO fully differential inverter cell and its loading. Four pMOS transistors 9, 11, 13, 15 have their drains tied to the voltage Vdd. The gates of the transistors 9 and 15 are both tied to a voltage Vb 19. The voltabe Vb 19 is generated from a voltage Vbn 18 through a replica bias. Here, Vbn 18 is the control voltage for controlling the current Icontrol. The gate of transistor 11 is tied to the sources of the transistors 9 and 11 as well as to the output 8 of the differential outputs 10 and 8. The gate of transistor 13 is similarly tied to the sources of the transistors 13 and 15 as well as to the output 8 of the differential outputs 10 and 8. A capacitor 16 is connected between the differential outputs 10 and 8. This capacitor actually reduces the output frequency of the CCO 7, however, it is necessary for improving the jitter performance.\nThe nMOS transistors 12, 14 have gates supplied by current supply inputs 2 and 3 which are connected to the output of the previous stage of the inverter cell as illustrated in FIG. 9. The sources of the transistors 12, 14 are connected to the sources of the transistors 9, 11 and 13, 15, respectively. The source of the transistor 12 also leads to the differential output 8. Connected to the drains of the transistors 12, 14 is the source of another nMOS transistor 16 having its gate supplied by a voltage 18. The transistor 16 has its drain grounded.\nFIG. 6 illustrates the CCO gain of the prior art circuit 7 of FIG. 1. Control current (in amps) is plotted along the x-axis while frequency (in Hertz) is plotted along the y-axis. There are separate curves for different design process corners and temperatures. The curves represent the SSH (slow-slow, high temperature), normal and FFL (fast-fast, low temperature) conditions. The curves, especially for the SSH condition tend to flatten when the control current becomes large. This is because the transistors 9, 15 enter the velocity saturation and their gm value does not continue to increase with control current.\nAn example of a prior art CCO design providing temperature variation compensation is presented in the paper entitled, “A 622-MHz Interpolating Ring VCO with Temperature Compensation and Jitter Analysis”, by Wing-Hong Chan, published in the IEEE International Symposium on Circuits and Systems, Jun. 9-12, 1997, Hong Kong. However, the method this paper can only provide compensation at one fixed frequency and cannot compensate for process variation. In addition, it requires many additional circuits resulting in greater power consumption, size and cost.\nAnother example of a prior art CCO design is presented in “Low-Jitter Process-Independent DLL and PLL Based on Self-Bias Techniques', IEEE J. Solid-State Circuits, vol. 31, No. 11, November 1996 by John G. Maneatis. However, this prior art CCO does not sufficiently extend the linear region of the CCO gain or minimize the process and temperature sensitivity.\nIt would therefore be desirable to provide a CCO with extended linear gain over a broad tuning range, greater stability, reduced size and power consumption and reduced sensitivity to process and temperature variation. Additionally, it would be desirable to provide a CCO with these features while maintaining a good power supply rejection ratio (PSRR)."}
{"text": "My invention relates to shaving mugs. The primary objectives are to eliminate the adverse characteristics of the currently existing mugs and provide a person with a reuseable shaving mug designed specifically to hold a typical shaving soap cake in a manner that is truly convenient to the user, even in travel environments. The adverse characteristics of the currently existing shaving mugs are:\na. The soap cake easily dislodges allowing the soap cake to flop around during use or to fall into the sink during clean-up. PA1 b. A lid is not provided or is not secured, allowing other items in luggage to become soiled when traveling. PA1 c. Breaks easily if dropped during use or luggage handling. PA1 a. Prevent the soap cake from dislodging during use, . . . cleaning, . . . handling, . . . or storage. PA1 b. Provide a space to contain the lather during use. PA1 c. Permit the soap cake to dry quickly between use. PA1 d. Protect the soap cake from contamination between use, especially during travel. PA1 e. Prevent soap and water from soiling clothing and other items in luggage during travel. PA1 f. Not break when dropped from normal sink heights or during typical rough handling of luggage."}
{"text": "Diffractive Optical Elements (DOEs) are now widely used in technical optics because they are planar elements which can easily be mounted, packaged, because they can be made by means of high productivity batch planar technologies, and also because they can perform complex optical functions which refractive light processing techniques can not, or can hardly achieve by simple and low cost means.\nOne of the hurdles in diffractive optics technologies is the difficulty of obtaining high diffraction efficiency. One of the widely used solutions is to decompose the desired analogue groove profile in a number of discrete levels in a staircase form. This is a costly multilevel manufacturing process which is presently limited to elements of small angular aperture. Another solution is to rely upon a gray scale process technology capable of photolithographically and physically transferring the desired analogue profile with fidelity in one single technological step. This solution is not available yet as a manufacturing technology. Yet another solution is to replicate the mould of an analogue surface relief generated by means of highly resolving writing means such as an electron beam pattern generator. This technology is only suitable for manufacturing large volumes of identical elements and suffers from the thermal instability and ageing of the replicable material.\nThere is therefore the need for a technical solution providing high diffraction efficiency gratings, diffractive optical elements, holograms, even in the presence of more than one propagating diffraction order without having to resort to a multilevel technology to generate analogue groove profiles.\nAn optical diffraction device is described in U.S. Pat. No. 6,219,478 B1 which discloses a reflective diffraction grating where the diffraction event provides large, possibly 100% diffraction efficiency for the sole −1st propagating diffraction orders in the incidence medium in a direction which is not parallel to the incident wave. This document U.S. Pat. No. 6,219,478 B1 discloses the conditions for the incident beam field to accumulate into a leaky mode of a layer placed on top of the mirror. A leaky mode is a transverse field resonance leaking into the cover medium. A diffraction grating written on, or within the layer acts as a tap regulating the rate of field accumulation so that the accumulated field leaking into the incident medium damps or cancels out the reflection in the direction of the Fresnel reflection by destructive interference between the wave directly reflected from the top surface of the layer and the leaking wave which is accumulated in the leaky mode. Consequently, and provided the incidence configuration is not autocollimation, the optical energy has nowhere else to propagate but to be directed along the −1st diffraction order of the grating.\nAlthough U.S. Pat. No. 6,219,478 B1 discloses the means to possibly achieve up to 100% diffraction efficiency in a grating even when the incidence angle θc (the incidence angle θc is defined from the normal to the general plane in which the grating extends) is relatively small, there are structures where the refractive index difference between the leaky mode propagating layer and the cover medium is too small, and/or there are incidence configurations where the incidence angle is too small, to provide sufficient field accumulation in the leaky mode to permit the cancellation of the reflection, therefore to give rise to 100% diffraction efficiency. In cases where the reflection can nevertheless be cancelled, deep grooves are required which implies that the leaky mode resonance is spectrally and angularly broad. According to U.S. Pat. No. 6,219,478 B1, it is under substantially grazing incidence and/or in the presence of a large index contrast between the leaky mode propagating layer and the cover medium that the reflection coefficient of the top boundary of the dielectric layer is large, i.e., that close to 100% diffraction efficiency can be achieved.\nIt would therefore be advantageous to achieve the cancellation of the reflection on the leaky mode propagating layer in all cases where the incident or the diffracted beam angle is small and in cases where the refractive index contrast between leaky mode propagating layer and cover is small (as for instance in the case of holograms and most visual diffractive structures which are often coated by a protective layer; when such a coating is applied on top of the leaky mode propagating layer, it is even impossible to obtain large incidence angle at the top surface of this layer since it would correspond to total reflection at the air-coating interface), and to achieve high, possibly 100% −1st order diffraction efficiency by means of a relatively weak corrugation or index modulation (besides, it is not always desired or possible to fabricate a deep corrugation or to cause a large index modulation in the layer propagating the leaky mode).\nThe documents of the scientific literature dealing with high diffraction efficiency of the −1st order of a reflection grating usually consider the diffraction configuration of the −1st order Littrow incidence where the diffracted beam is diffracted back in the direction of the incident beam. The Littrow incidence condition for the −1st order at vacuum wavelength λ from an incidence medium of refractive index nc on a periodic grating of period Λ is characterized by the Littrow angle θL such that sin(θL)=λ/(2Λnc). This incidence condition is also currently called the autocollimation configuration since the reflected −1st order diffracted wave propagates back parallel to the incident wave. The off-Littrow incidence configuration will hereafter refer to configurations where the angle of incidence θc essentially differs from θL, i.e., where the −1st order reflected diffracted wave is not parallel to the incident wave. Apart from a major functional difference, there is a fundamental difference between the autocollimation diffraction configuration and the off-Littrow configuration. Unlike in the off-Littrow diffraction configuration, the autocollimation configuration is known to always permit 100% diffraction efficiency (provided the layer average thickness is larger than a minimum thickness and provided the sole 0th and −1st orders can propagate) as calculated for instance in the case of high efficiency femtosecond pulse compression gratings in document M. D. Perry, R. D. Boyd, J. A. Britten, D. Decker, B. W. Shore, C. Shannon, E. Shults, “High-efficiency multilayer dielectric diffraction gratings”, Opt. Lett. 20 No 8, 940-942 (1995) and as analyzed by document H. Wei, L. Li, “All-dielectric reflection gratings: a study of the physical mechanism for achieving high efficiency”, Appl. Opt., Vol. 42 No 31, 6255-6260 (2003) which explains the mechanism of high efficiency and limits itself to the Littrow case. So does the document by V. A. Sychugov, B. A. Usievich, K. E. Zinov'ev, O. Parriaux, “Autocollimation diffraction gratings based on waveguides with leakage mode”, Quantum Electronics, 30(12) 1094-1098 (2000) which reports on the use of a semi-reflective structure on top of a dielectric layer with the objective of obtaining the achievable 100% diffraction efficiency in the −1st order Littrow configuration by means of a grating of a smaller corrugation depth. The semi-reflective structure consists of quarter wave layers at the wavelength, and at the incidence and diffraction direction of the autocollimation configuration.\nThe reason for such specific feature of the Littrow configuration is that the field of the diffracted wave in the mirrored corrugated structure is the same as that of the incident wave. For instance, the mirror ensuring the reflection of the incident wave reflects the diffracted wave identically. Similarly, in the structure dealt with by the above mentioned paper by V. A. Sychugov et al, the semi-reflective structure ensuring some degree of field concentration in the leaky mode inherently ensures the same degree of field concentration for the diffracted leaky mode. The autocollimation configuration is analogous to the reflection from a quarter wave multilayer mirror or from a fibre Bragg grating mirror with the specificity that the two waves participating in the −1st order reflection (the −1st order directed into the cover and the −1st order directed to the mirror, then reflected into the cover) have to interfere essentially constructively.\nBecause its symmetry, the autocollimation configuration implies first order coupling between the two counterpropagating leaky modes if the leaky mode propagation condition is satisfied. It is to be noted that the autocollimation configuration permits to obtain 100% diffraction efficiency even if the leaky mode propagation condition is not satisfied.\nIt is an object of the invention to provide an optical diffraction device having high and possibly 100% diffraction efficiency for the −1st diffraction order in a diffraction configuration where neither the incident beam nor the diffracted beam are grazing, where the diffracted beam is not parallel to the incident beam (i.e., outside the Littrow configuration), by means of a diffractive element or structure having relatively shallow depth or/and weak index modulation."}
{"text": "The present invention relates to speakers and in particular to speakers of the moving coil type of relatively small diameter for use in stereo headsets and the like.\nHeretofore it has been common practice to construct the diaphragm assembly for a moving coil type of headset speaker of three discrete parts. These parts comprise a low frequency diaphragm a high frequency dome and a voice coil form about which the voice coil is wound. The three parts were cemented together to form the completed diaphragm assembly. In this prior art type of construction, two distinct gluing operations were required. That is, the high frequency dome was separated from the voice coil form by adhesive and the low frequency diaphragm was also separated from the voice coil form by an adhesive layer.\nIt is common knowledge that in a mechanical acoustic driving system, such as a loudspeaker, a phase shift occurs at each change in the density of the material from which the driver is formed. This phase shift results is distortions at specific frequencies.\nIt is also well known that it is desirable to keep the vibrating system (i.e., the diaphragm-dome assembly) as light-weight as possible in order to achieve the lowest possible free field resonance frequency. With conventional headset two inch speakers now available, this frequency is approximately 450 Hz."}
{"text": "The present disclosure relates to a measurement system for the measurement of dry powder based agents.\nIn order to certify a dry powder fire suppression system onboard a vehicle such as an aircraft, the agent is discharged into the protected volume and an analyzer simultaneously records the amount of fire extinguishing agent in various zones of the protected volume. The amount of agent must be above some predetermined level which has been established sufficient to extinguish all possible fires for some period of time simultaneous in all zones.\nThe analyzer must be calibrated and traceable such that analyzer output proves the dry powder fire suppression system is capable of extinguishing any fire within the protected space. No known systems are capable of both measuring aerosol cloud fire extinguishing agent concentrations and being calibrated so as to measure the agent concentration for an aircraft dry powder fire suppression system certification test."}
{"text": "The present invention relates to a single-output resonant switched power converter, in which the transfer of energy between the input and the output thereof is governed by a switching element located on the primary side of the resonant switched power converter.\nThe secondary side of the converter includes a resonant circuit that carries out the demagnetisation of the core of a transformer, which electrically isolates the input from the output of the resonant switched power converter.\nThe resonant switched power converter is of special application, but not exclusively, in distributed power supply systems in which switched power converters with low voltage and power consumption are required, since they are mounted on printed circuit boards as one more electronic component thereof, and consequently a high integration density and reduced dimensions are required.\nA switched power supply converter is known from the U.S. Pat. No. 5,886,881 granted to Xia et al., incorporated in the present patent application by reference. It describes a single-output forward converter that comprises on the primary side a serial combination of a primary winding of a transformer and a first switching element, which is connected to an input voltage supplied by a direct current (DC) power source.\nOn the secondary side, a secondary winding of the transformer is connected in cascade to a self-excited synchronous rectifier and to a filter that applies an output voltage to a load. The output of the filter constitutes the output of the switched power converter.\nA serial combination consisting of a capacitor and a second switching element is connected in parallel with a third switching element, which forms the rectifier branch of the self-excited synchronous rectifier.\nWhen the first switching element is in a conducting state, the input voltage is applied to the input of the self-excited synchronous rectifier through the transformer. During this period there is direct energy transfer between the input and the output.\nWhen the first switching element is in a non-conducting state, the second switching element is in conduction, a magnetising current flowing from the secondary winding to the capacitor, which is charged to a voltage proportional to the input voltage. The result being that the voltage on the secondary winding is maintained constant, or clamped, and the core of the transformer is demagnetised, or reset.\nA drawback of the forward switched power converter is that the non-conducting period of the first switching element is constant and, as a consequence, the demagnetisation of the core of the transformer does not take place during the entire non-conducting period of the first switching element, the result being that in the event of major variations of the input voltage, undesirable dead times appear in the switching of the self-excited synchronous rectifier.\nFor this reason, it is necessary to develop a switched forward power converter which accepts a broad range of input voltages and guarantees, for all of them, that the demagnetisation of the core of the transformer takes place throughout the entire non-conducting period of the first switching element, in such a manner that in the operation of the self-excited synchronous rectifier no dead times appear over the whole range of input voltages. All this leading to an enhancement of the forward converter performance, while maintaining both simplicity of design and of operation of the forward converter.\nTo overcome the problems described above, a resonant switched power converter is proposed that is capable of making the end of the resonant period coincide with the start of the conducting period of a first switching element, thereby avoiding undesirable dead times during the operation of a self-excited synchronous rectifier.\nTo achieve said objective a first capacitor is added which has the characteristic that its capacitance is a function of the duty cycle of the first switching element. To achieve said objective the first capacitor is connected in parallel with a second winding of a transformer.\nThe variation of the value of the capacitance of the first capacitor is carried out by adding in parallel with said first capacitor a serial combination of a second capacitor and a second switching element. The value of the capacitance varies between a first value and a second value.\nBy controlling the value of the capacitance of the first capacitor, the end of the resonant period is made to coincide with the end of the nonconducting period of the first switching element."}
{"text": "A wide variety of mental and physical processes are controlled or influenced by neural activity in particular regions of the brain. For example, various physical or cognitive functions are directed or affected by neural activity within the sensory or motor cortices. Across most individuals, particular areas of the brain appear to have distinct functions. In the majority of people, for example, the areas of the occipital lobes relate to vision; the regions of the left interior frontal lobes relate to language; portions of the cerebral cortex appear to be consistently involved with conscious awareness, memory, and intellect; and particular regions of the cerebral cortex, the basal ganglia, the thalamus, and the motor cortex cooperatively interact to facilitate motor function control.\nMany problems or abnormalities with body functions can be caused by damage, disease, and/or disorders in the brain. For example, Parkinson's Disease (PD) is related to the degeneration or death of dopamine producing neurons in the substantia nigra region of the basal ganglia in the brain. Dopamine is a neurotransmitter that transmits signals between areas of the brain. As the neurons in the substantia nigra deteriorate, the reduction in dopamine causes abnormal neural activity that results in a chronic, progressive deterioration of motor function control. Conservative estimates indicate that PD may affect more than one million individuals in the United States alone.\nPD patients typically exhibit one or more of four primary symptoms. One primary symptom is a tremor in an extremity (e.g., a hand) that occurs while the extremity is at rest. Other primary symptoms include a generalized slowness of movement (bradykinesia); increased muscle rigidity or stiffness (rigidity); and gait or balance problems (postural dysfunction). In addition to or in lieu of these primary symptoms, PD patients may exhibit secondary symptoms including: difficulty initiating or resuming movements; loss of fine motor skills; lack of arm swing on the affected side of the body while walking; foot drag on the affected side of the body; decreased facial expression; voice and/or speech changes; cognitive disorders; feelings of depression or anxiety; and/or other symptoms.\nEffectively treating PD or other movement disorders related to neurological conditions can be very difficult. Current treatments for PD symptoms include drugs, ablative surgical intervention, and/or neural stimulation. Drug treatments or therapies may involve, for example, the administration of a dopamine precursor that is converted to dopamine within the central nervous system (i.e., Levodopa (L-dopa)). Other types of drug therapies are also available. Unfortunately, drug therapies frequently become less effective or ineffective over time for an undesirably large patient population. A PD patient may require multiple drugs in combination to extend the time period of efficacy of drug therapies. Drug treatments additionally have a significant likelihood of inducing undesirable physical side effects; motor function complications such as uncontrollable involuntary movements (dyskinesias) are a particularly common side effect. Furthermore, drug treatments may induce undesirable cognitive side effects such as confusion and/or hallucinations.\nAblative surgical intervention for PD typically involves the destruction of one or more neural structures within the basal ganglia or thalamus that have become overactive because of the lack of dopamine. Unfortunately, such neural structures reside deep within the brain, and hence ablative surgical intervention is a very time consuming and highly invasive procedure. Potential complications associated with the procedure include risk of hemorrhage, stroke, and/or paralysis. Moreover, because PD is a progressive disease, multiple deep brain surgeries may be required as symptoms progressively worsen over time. Although ablative surgical intervention may improve a PD patient's motor function, it is not likely to completely restore normal motor function. Furthermore, since ablative surgical intervention permanently destroys neural tissue, the effects of such intervention cannot be readily adjusted or “fine tuned” over time.\nNeural stimulation treatments have shown promising results for reducing some of the symptoms associated with PD. Neural activity is governed by electrical impulses or “action potentials” generated in and propagated by neurons. While in a quiescent state, a neuron is negatively polarized and exhibits a resting membrane potential that is typically between −70 and −60 mV. Through chemical connections known as synapses, any given neuron receives excitatory and inhibitory input signals or stimuli from other neurons. A neuron integrates the excitatory and inhibitory input signals it receives, and generates or fires a series of action potentials in the event that the integration exceeds a threshold potential. A neural firing threshold, for example, may be approximately −55 mV. Action potentials propagate to the neuron's synapses and are then conveyed to other synaptically connected neurons.\nNeural activity in the brain can be influenced by neural stimulation, which involves the application of electrical and/or magnetic stimuli to one or more target neural populations within a patient using a waveform generator or other type of device. Various neural functions can thus be promoted or disrupted by applying an electrical current to one or more regions of the brain. As a result, researchers have attempted to treat certain neurological conditions, including PD, using electrical or magnetic stimulation signals to control or affect brain functions.\nDeep Brain Stimulation (DBS) is a stimulation therapy that has been used as an alternative to drug treatments and ablative surgical therapies. In DBS, one or more electrodes are surgically implanted into the brain proximate to deep brain or subcortical neural structures. For treating PD or other movement disorders, the electrodes are positioned in or proximate to the ventrointermediate nucleus of the thalamus; basal ganglia structures such as the globus pallidus internalis (GPi); or the Subthalamic Nucleus (STN). The location of the stimulation site for the electrodes depends upon the symptoms that a patient exhibits and the severity of the symptoms.\nIn a typical DBS system, a pulse generator delivers a continuous or essentially continuous electrical stimulation signal having a pulse repetition frequency of approximately 100 Hz to each of two deep brain electrodes. The electrodes are may be positioned bilaterally on the left and right sides of the brain relative to particular neural structures such as those indicated above. U.S. Pat. No. 5,883,709 discloses one conventional DBS system for treating movement disorders.\nAlthough DBS therapies may significantly reduce one or more PD symptoms, particularly when combined with drug treatments, they are highly invasive procedures. In general, configuring a DBS system to properly function within a patient requires two time consuming, highly invasive surgical procedures for implanting the DBS electrodes. Each such surgical procedure has essentially the same risks as those described above for ablative surgical intervention. Moreover, DBS may not provide relief from some movement disorders.\nMotor Cortex Stimulation (MCS) is another type of brain stimulation treatment that has been proposed for treating Parkinson's Disease. MCS involves the application of stimulation signals to the motor cortex of a patient. One MCS system includes a pulse generator connected to a strip electrode that is surgically implanted over a portion of only the motor cortex (precentral gyrus). The use of MCS to treat PD symptoms is described in Canavero, Sergio, Extradural Motor Cortex Stimulation for Advanced Parkinson's Disease: Case Report, Movement Disorders (Vol. 15, No. 1, 2000).\nBecause MCS involves the application of stimulation signals to surface regions of the brain rather than deep neural structures, electrode implantation procedures for MCS are significantly less invasive and time consuming than those for DBS. As a result, MCS may be a safer and simpler alternative to DBS for treating PD symptoms. Present MCS techniques, however, fail to address or adequately consider a variety of factors that may enhance or optimize the extent to which a patient experiences short term and/or long term relief from PD symptoms."}
{"text": "Television broadcasting has evolved from basic analogue terrestrial broadcast television to complex digital systems. Wireless communication techniques are central to the development of the complex digital systems. There exists several wideband digital communication techniques depending on a broadcasting method used. For example, direct sequence spread spectrum (DSSS) and orthogonal frequency-division multiplexing (OFDM) are one of the latest schemes in wideband digital communication systems. OFDM is a method of encoding digital data on multiple carrier frequencies and is used in applications such as digital television and audio broadcasting, Digital Subscriber Line (DSL) internet access, wireless networks, power line networks, and 4G mobile communications. OFDM has been selected as the wireless technique for the current generation of terrestrial television broadcast standards such as DVB-T2 and emerging standards such as ATSC 3.0.\nA broadcasting standard may allow many modes of operations to be determined by the broadcaster. Thus, signaling parameters sire needed b the receiver to decode efficiently, and correctly the received data. As recognized by the present inventors, there is a need to transmit signaling parameters.\nThe foregoing “Background” description is for the purpose of generally presenting the context of the disclosure. Work of the inventors, to the extent it is described in this background section, as well as aspects of the description which may not otherwise qualify as prior art at the time of filing, are neither expressly or impliedly admitted as prior art against the present disclosure. The foregoing paragraphs have been provided by way of general introduction, and are not intended to limit the scope of the following claims. The described embodiments, together with further advantages, will be best understood by reference to the following, detailed description taken in conjunction with the accompanying drawings."}
{"text": "Field of the Invention\nThe present invention is directed generally to fluoroscopic systems and more particularly to a large field, low intensity, radiation imaging system.\nU.S. Pat. No. 4,142,101 to Yin discloses a low intensity x-ray and gamma-ray imaging device. The imaging device includes a phosphor screen for converting an x-ray image into a visible light image. The resulting light is transmitted via an input fiber optic plate to a photocathode where it is converted to electrons. The electrons are accelerated through a vacuum to a microchannel plate image intensifier which multiplies the electrons. The resulting electrons are then converted to light by a second phosphor screen and the light is thereafter directed via an output fiber optic plate to a viewer. The low intensity x-ray and gamma-ray imaging device disclosed in Yin is also described in Yin et al., \"The Lixiscope\", Nasa Technical Memorandum 79634, September 1978.\nThe device disclosed in U.S. Pat. No. 4,142,101 has become known as a Lixiscope which stands for low intensity x-ray imaging scope. Lixiscopes are commercially available from HealthMate, Inc., 3000 Dundee Road, Northbrook, Ill. 60062. Lixiscopes provide substantial advantages over traditional x-ray imaging systems Lixiscopes are small, lightweight, and portable. In addition, because of the microchannel plate image intensifier, a very low level source of x-rays or gamma radiation can be used thus reducing the overall cost of the system and decreasing the exposure of the patient. Unfortunately, Lixiscopes are limited to small fields of view, typically about two inches. Thus, Lixiscopes do not compete with large format medical fluoroscopic devices.\nLarge field x-ray images can be produced by an x-ray image amplifier such as that disclosed in U.S. Pat. No. 2,681,868. Such x-ray image amplifiers allow a large field of view but require sizable x-ray exposure of the patient. Also, the x-ray generator must be sufficiently large to produce the required x-rays which increases the cost of the overall system. Despite the exposure of the patient to a large dosage of x-rays, it is oftentimes necessary to darken the background such that the x-ray image can be properly viewed. Thus, it is desirable to combine the large field of view of traditional fluoroscopic devices with the low radiation exposure, high gain, and low cost of the Lixiscope.\nOne attempt at combining the large field of view of traditional fluorscopic devices with the advantages of the Lixiscope is disclosed in U.S. patent application Ser. No. 738,616 for a Large Screen Microchannel Plate Radiation Imaging System filed May 28, 1985, and assigned to the same assignee as the present invention. In U.S. patent application Ser. No. 738,616 a method and apparatus are disclosed in which invisible radiation from a radiation source is passed through an object to form a large field invisible radiation image of the object. The invisible radiation image is converted to a visible light image. The visible light image is reduced in size using lenses, mirrors, fiber optic tapers, or some combination thereof. The reduced visible light image is then intensified and the intensified image may be projected for viewing.\nWhen a fiber optic taper is used for reducing the size of the visible light image, optical fibers are typically subjected to an additional manufacturing step which draws the fibers such that the input ends of the optical fibers have a greater diameter than the output ends of the optical fibers. Although this does effect size reduction, the drawing of the optical fibers results in optical fibers having lower numerical apertures i.e. optical fibers which transmit less light. Also, because of radiation scattered by the object being viewed, the image of the object may be blurred. This blurring may be reduced by using a grid. However, the grid, in addition to eliminating the scattered radiation, also prevents some of the desired radiation from being transmitted which may result in degradation of the image and loss of brightness. Finally, if the grid is not oscillated during exposure of the object to the radiation, the grid pattern will be superimposed on the resulting image.\nA different approach to providing a large field, low intensity, radiation imaging apparatus centers around the reduction of the electron image corresponding to the radiation image of interest rather than the reduction of the visible light image as set forth in the above-identified patent application. An apparatus directed to the approach of reducing the size of the electron image is disclosed in U.S. patent application Ser. No. 849,907 for a Large Field, Low Intensity, Radiation Imaging System filed Apr. 9, 1986, and assigned to the same assignee as the present invention. The apparatus disclosed in that patent application converts a radiation image to an electron image. Circuitry is provided for reducing the size and increasing the intensity of the electron image. Additional circuitry outputs the increased intensity electron image.\nThe apparatus of Ser. No. 849,907 may suffer from degradation of the output image due to scattered radiation. Although the scattered radiation may be substantially eliminated by the use of a grid, the grid will also interfere with certain of the radiation which is a portion of the image to be viewed. Thus, the grid will result in signal degradation and grid lines will be superimposed on the image as in U.S. patent application Ser. No. 738,616. Oscillation of the grid washes the grid lines from the image but also degrades the image and results in loss of brightness.\nIt therefore remains desirable to provide a large field, low intensity, radiation imaging system which takes advantage of the improvements provided by using a grid but which does not suffer from any degradation of output or loss of brightness as a result of using the grid."}
{"text": "Efficient orderly resource management of diagnostic tests in human disease is a resistant problem. Many diagnostic tests applied on human patient populations carry an element of risk. The decision matrix by which diagnosticians reduce symptoms to specific diagnostic tests varies due to differences in medical centers' diagnostic capacities. Health Maintenance Organization (HMOs), Preferred Physician Providers Organizations (PPOs), and others' methods introduced to reduce healthcare resource consumption by fiscal means have primarily failed as measured by the gross national product. What is needed in the medical arts is an improved method of disease diagnosis that is more effective and efficient than methods currently employed by diagnostic centers."}
{"text": "Rubber polymers are used in the treatment of millions of patients each day. Bacterial and protein adherence to polymer materials, implanted or inserted into the body, are an early step in blood clots and biofilm formation, and typically result in infections, complications from the resultant infections, longer hospital stays, and perhaps death. Due to protein adsorption and subsequent bacterial adherence, biofilms are formed at interfaces between solid substrates and liquids. The prevention of bacterial and protein adherence and biofilm development on biomedical materials continues to be a major challenge in healthcare. Catheter-associated urinary tract infections (CAUTI) alone are the most common health-care-associated infection worldwide. Accordingly, there exists a significant need for new anti-infective medical grade polymer biomaterials."}
{"text": "1. Field of the Invention\nThis invention relates to a trolley in particular to an accident and emergency trolley.\n2. Description of Related Art\nIt is known to provide accident and emergency trolleys designed to be raised and lowered, tilted longitudinally, but such conventional trolleys suffer from the disadvantages of being difficult to operate and manouver and the patient having to be transferred to another surface for X-rays or imaging.\nAccordingly, the present invention provides an accident and emergency trolley for supporting and transporting a patient comprising a surface overlying an upper frame, the upper frame supported on a base frame so that it can be raised and lowered relative to the base frame, the upper frame being mounted at its periphery on opposed radius arms connected to respective first and second actuators, the actuators operated by foot operated pedals, said pedals is disposed on one or either side of the trolley such that operation of any single pedal enables the trolley to be raised at either end or both ends simultaneously. Advantageously, the opposed radius arms provide an uninterrupted window under the trolley allowing for X-ray/imaging. Further their operation is much simplified only requiring a single operation compared to the complex arrangement of pedals on conventional trolleys in order to achieve tilt. Preferably, the pedals for all functions are advantageously located at each side of the trolley allowing for easy access and avoids the operators having to travel and operate pedals located around the trolley for individual functions.\nPreferably, the pedals may only be operated with a downward force, thereby providing easier operation. Previous trolleys have pedals which need to be pushed up to operate actuators for certain tilt functions.\nPreferably, the trolley comprises a castor mounted centrally on the base frame by resilient means and arranged to engage in either of two positions 180xc2x0 apart so as to provide directional tracking in line with longitudinal axis of the trolley. Therefore, if pushed in a longitudinal direction the castor automatically latches and provides improved steering and maneuverability.\nAdvantageously, the trolley surface may be sectional to provide at least a backrest movable from a rest position to an inclined position, by means of actuators operable by a single push action, for ease of operation.\nPreferably, the trolley additionally comprises safety side rail assemblies each comprising rail arms pivotally mounted on a support member suspended from the upper frame by links, the rail arms movable from a horizontal stowed position to a vertical in-use position, the support member being linked to the upper frame such that it is stowable under its own weight, the arms being locked in their in-use position.\nThe side rail assemblies advantageously provide for safety of the patient but can be stowed away under the upper frame so that there is no transfer gap when transferring a patient from the trolley to an operating table, bed or similar."}
{"text": "The present invention relates to a method and a device for making shaped articles from thermoplastic plastics, particularly non-cleaned, non-sorted plastic refuse.\nWith methods and devices of the type under consideration plastic refuse is plasticized in a roller extruder and fed into molds in which pressure is applied to the plastic filling the molds which are afterwards cooled and from which shaped objects are thereafter removed.\nMethods and devices for a repeated use of plastic refuse or plastic trash have been not, however, known until now.\nIt is known that plastic refuse which was cleaned, sorted and re-granulated was treated together with a new plastic to produce plastic articles. Since the cleaning and sorting involves a manual work and the plastic refuse or trash is a garbage-like or trash-like non-homogeneous substrate no working forces are normally available in the economy to do this kind of work; and when such forces are found this kind of work is too expensive. A remarkable portion of new plastic does not solve the problem because the product consisted of such re-granulated material cannot be successfully marketed due to its quality and is a so-called second choice. It is particularly difficult to use this re-granulate due to high costs of production.\nRe-granulates are as a rule plasticized for processing in worm-type extruders.\nThe solution of the problem has been found in the field of elimination of industry and household thermoplastic refuse and trash in the utilization of roller extruders disclosed, for example, in DE-PS No. 27 00 846. The roller extruder is so robust that it not only is able to process non-cleaned, non-sorted thermoplastic trash or refuse but also rock and metal particles included in the plastic trash or refuse so that the plasticized thermoplastic plastic remains homogeneous and has no quality deficiencies. The plasticized material pushed from the nozzle of the roller extruder, however, has been until now, calibrated into pieces of suitable size by hand, cut and fed into respective molds in which the plastic was pressed by suitable devices. The objects formed of plastic normally had the weight of less than one kilogram.\nDespite the homogenizing and high pressures used for producing shaped articles such articles are of pure quality and no market exists for such articles. Such articles can be relatively expensive when there is a chance of amortization of initial costs; they should however not be expensive on the market where a consumer is found.\nThe thermoplastic refuse material plasticized in the roller extruder can result in the material which is acid-resistant, weather-resistant and resistant to insects and further damaging factors and which would have high rigidity.\nIn order to solve the problem of increasing amounts of collected plastic refuse or trash the elimination of which by burning is costly and problematic it has been suggested to introduce this refuse into the economical circuit as a worthwhile crude product which would find a consumer. It has been found that such product pallets should consist of articles which are large and heavy not only because a great deal of plastic refuse is consumed but also because heavy articles, for example having weight of 100 or many hundred kilograms are not of a third choice as it concerns their quality, but their properties, which affect quality of small products, can be utilized.\nThree groups of shaped objects were selected as preferable product pallets, namely universal plates, solid profiles and hollow bodies. The solid profiles can be utilized similarly to bricks for making palisades; they can form non-prismatic shaped table plates, warning signal pedestals or the like.\nIn the production of these objects which has been carried out by methods as described herein above, such objects are always too expensive to be on the market, and practically it has been impossible to find consumers for such objects."}
{"text": "Heretofore, there has been proposed a drum-shaped roller rotatable in forward and backward directions as a means for making input of character information etc.\nAn input means of this type is known from the disclosure in the Japanese Published Unexamined Patent Application No. 203985 of 1999.\nFIGS. 1 to 3 illustrate together a portable phone as an electronic apparatus using a rotating drum of the above type as a means for making entry of character information. FIG. 1 shows the external view of the portable phone, FIG. 2 is a circuit diagram of the portable phone in FIG. 1, and FIG. 3 shows a multi-functional switch used as a mobile communications terminal.\nAs shown in FIGS. 1 and 2, the portable phone includes an antenna 1, transmission/reception circuit 2, microphone 3, speaker 4, CPU (central processing unit) 5, ROM (read-only memory) 6 having stored therein programs etc. executed by the CPU 5, RAM (random-access memory) 7 having provisionally stored therein data required for execution of a program by the CPU 5, display drive circuit 8, and a display unit 9 driven by the display drive circuit 8 to display a variety of information. The display unit 9 uses a liquid crystal display (LCD) panel. The portable phone shown in FIG. 1 further includes operating keys 10 and a multi-functional switch 11 forming the information input device. The operating keys 10 and multi-functional switch 11 form together an operation unit 12.\nIn the portable phone having a body 13, the antenna 1 is provided at the right upper end of the body 13 as shown in FIG. 1. The antenna 1 is extensible upward. Also, in this portable phone, the speaker 4 is disposed nearly in the middle of the upper end of the apparatus body 13, the display unit 9 is disposed below the speaker 4, the operation unit 12 is disposed below the display unit 9, and the microphone 3 is disposed in the middle of the lower end of the apparatus body 13. The multi-functional switch 11 is disposed in the upper portion of the operation unit 12 and the operating keys 10 including ten keys are disposed below the multi-functional switch 11.\nFIG. 3 shows the multi-functional switch 11 in detail. As shown, it has a drum-shaped roller 11A. By rotating or pressing the roller 11A, information is entered. The multi-functional switch 11 is constructed and functions as will be described below:\nAs shown in FIG. 3, the multi-functional switch 11 includes a base 11D having a downward-extending fulcrum projection 11F formed along one longitudinal side thereof. The base 11D is installed to a printed circuit board 11E with the fulcrum projection 11F laid between them. The base 11D is normally kept generally parallel to the printed circuit board 11E by a forcing means (not shown). However, when the base is pressed from above against the force of the forcing means, it is pivoted downward about the fulcrum projection 11F.\nThe roller 11A is supported above the base 11D rotatably with the rotating shaft thereof being held horizontally. More specifically, the roller 11A is supported at one end of the rotating shaft thereof in a bearing provided in a support piece 11G and at the other end in a bearing provided in a rotation detector 11B. The rotation detector 11B supporting the other end of the rotating shaft detects a rotated extent and direction of the rotating shaft of the roller 11A. A switch 11C is provided on the printed circuit board 11E and under the base 11D. When the roller 11A is pressed from above, the base 11D is pivoted about the fulcrum projection 11F and presses the switch 11C which will thus be turned on. Thus, the multi-functional switch 11 makes entry of a rotated extent and direction of the roller 11A by rotating the latter and also enters information by pressing the roller 11A.\nThe operations of the portable phone shown in FIG. 1 will be explained below as to a phone directory search with reference to FIG. 4:\nFirst, the operating keys 10 of the operation unit 12 are operated in a predetermined manner to display an initial screen 9a set in the portable phone. A menu for the phone directory search is selected from the initial screen 9a and then set. Next, the roller 11A of the multi-functional switch 11 is rotated upward or downward in the plane of FIG. 3. At this time, there will be displayed on the display unit 9 information about initials of a destination's name like “a, i, u, e, o, ka, ki, . . . wa” (in the order of the Japanese 50-character “kana” syllabary) for example as options. The information appears scrolling, and the options are extracted from the RAM 7 correspondingly to a rotated extent of the roller 11A and displayed on the display unit 9. A reference number 9b indicates an example of the display on the screen of the display unit 9.\nWhen a desired initial is found on the screen, a cursor 15 is pointed to that initial by rotating the roller 11A. That is, the information is selected. A reference number 9c indicates an example of the display on the screen when the desired information has been selected. In this example, “sa” is selected. After completion of this selection, the roller 11A is pressed to set the selected initial “sa”. That is, when the roller 11A is pressed, the switch 11C is turned on, which is detected by the CPU 5. An option the cursor 15 is pointing to, information about the initial “sa” in this example, is set. With this setting, names (full name) having been entered under the initial (“sa” in this example) will be displayed. A reference number 9d indicates an example of the display. Then, the roller 11A is further rotated to move the cursor 15 to a desired names. A reference number 9e indicates an example of the selection. A name “Ikuo Sasaki” is selected in this example. With this setting, the roller 11A is pressed. The name “Ikuo Sasaki” is set in this example.\nWith the roller 11A further pressed in this condition, information corresponding to the selected name “Ikuo Sasaki” is read from the RAM 7 and displayed on the display unit 9, which example is indicated with a reference 9f. For selection of a desired option from the display screen, for example, sending a phone call or an electronic mail to “Ikuo Sasaki”, the cursor 15 is moved to the option and placed over the option. Thus, the desired option is selected. For example, the cursor 15 is moved to a phone number and placed over it, whereby sending of the phone call or electronic mail to a phone of that number will be selected.\nNote that some destinations have a plurality of phone numbers including his or her home phone number, office phone number, a phone number of his or her portable phone number and also an electronic mail (e-mail) address or the like, not any single phone number.\nThe cursor 15 is of course moved to a phone number displayed on the screen by rotating the roller 11A of the multi-functional switch 11. A reference number 9g in FIG. 4 indicates an example when such a destination is selected. It should be noted that selection of an electronic mail address in the information corresponding to a selected name will lead to selection of sending of an electronic mail to that address.\nWhen the roller 11A is pressed from above in the plane of FIG. 3 after the above selection, the selected option is set. Then, the set option is carried out, namely, a call is sent to the selected one of the phone numbers of Ikuo Sasaki in this example. A reference number 9h indicates an example of the display appearing on the display unit 9 at this time, in which a message “calling”, the name (“Ikuo Sasaki”) of the destination of the phone call and the destination phone number are displayed.\nWith the portable phone shown in FIG. 1, it is possible to select one of the plurality of options by rotating the roller 11A of the multi-functional switch 11 and set the selected option by pressing the roller 11A with the finger having rotated the roller 11A for the selection of the option. Thus, it is possible to rapidly send a call or electronic mail.\nThe mobile communications terminal such as the conventional portable phone is not advantageous in some respects as follows since it uses the multi-functional switch 11.\nFirstly, the use of the multi-functional switch 11 inhibits the mobile communications terminal from being designed thinner. There are demands for mobile communications terminals including the portable phone having a wider variety of functions, higher performance, easier operability and a thinner design. The aforementioned conventional portable phone can hardly meet such demands. More specifically, in the multi-functional switch 11, the roller 11A has to be rotatably supported, and there should be provided between the base 11D to support the roller 11A and printed circuit board 11E the switch 11C provided to set a selected option and which is turned on when the roller 11A is pressed. This construction itself makes it difficult to provide a thinner design of the portable phone. Namely, the thickness of the multi-functional switch 11 depends primarily upon the diameter of the roller 11A. A smaller diameter of the roller 11A will lead to a thinner design of the multi-functional switch 11. For assuring a large scrolling by one rotated extent of the roller 11A, however, the diameter of the latter has to be correspondingly large. If the diameter of the roller 11A is too small, the roller 11A itself will loose its easy operability.\nThe above mechanism of the roller 11A will be described in detail below:\nThe circumferential turn of the roller 11A, required for a scrolling of one line through the screen of the display unit 9, has to be larger than a predetermined value and may not be limitlessly smaller than the predetermined value. When the diameter of the roller 11A is smaller, the rotation angle of the latter, necessary for the cursor movement for one line, has to be larger correspondingly (for the reduction of the diameter). The smaller the diameter of the roller 11A, the smaller the maximum number of pulses the roller 11A can generate per turn becomes. Therefore, for movement of the cursor over many lines, the roller 11A should be rotated very much by the finger, which will lead to a longer time for selection of a desired option and cause the operation of the roller 11A to be more annoying.\nOn this account, it is assumed that the scrolling attained by one turn of the roller 11A of a small diameter is increased. In this case, however, the circumferential turn of the roller 11A for movement of the cursor 15 for one line will be smaller, which makes it difficult to accurately point the cursor 15 to a desired line. The cursor 15 will move and return excessively, which will lead to a poorer operability of the roller 11A.\nAs will be known from the foregoing, the multi-functional switch 11 is limited from being designed thinner and it will possibly be impossible to meet the future demand for a thinner design of the mobile communications terminal or the like such as the portable phone which uses the multi-functional switch 11.\nSecondly, since with the multi-functional switch 11, an option can be selected by rotating the roller 11A and the selected option can be set by pressing the roller 11A after the selection, so the selecting and setting operations cannot be distinguished from each other. The users will not possibly be able to make a distinction between the selecting and setting operations. Namely, the following will be possible. A user intending to rotate the roller 11A will eventually press the latter because he or she has rotated the roller 11A with an excessive force, and thus he will set an undesired option. On the other hand, a user intending to press the roller 11A in order to set a selected option will rotate the latter because he or she has applied a force to the roller 11A in a wrong direction, and thus he will select a wrong or undesired option."}
{"text": "Light emitting diodes (LEDs) emit light when electrical current flows through them and emit no light when no current flows through them. LEDs require control of the current flowing through them to obtain desired brightness. One way of controlling this current flow, and brightness, is by alternately providing current flow through a LED and stopping current flow through the LED in response to a pulse width modulated (PWM) control signal. The percentage of the alternating current flow can be described as a duty cycle of the current flow.\nA fifty percent (50%) duty cycle describes a current flow through the LED 50% of the time and no current flow through the LED the other 50% of the time. A 10% duty cycle describes a current flow through the LED 10% of the time and no current flow 90% of the time. A 90% duty cycle describes a current flow through the LED 90% of the time and no current flow 10% of the time. The PWM signal can have a duty cycle to control the current flow and the brightness of the LED.\nAlthough the controlling pulse width modulated signal may have a well-defined 50% duty cycle of on and off times, unequal capacitances in the circuits implementing the flow of current through the LEDs may cause unequal turn on and turn off times of the current, and produce other than a 50% duty cycle of the LEDs.\nWhile some integrated circuit LED controllers contain an integrated power transistor to drive external power to the LEDs, some applications use an integrated circuit to drive an external power transistor to drive external power to the LEDs.\nIn FIG. 1, circuit 10 provides a series connection between a power lead 12 supplying a voltage such as from a battery, Vbat, and a field ground lead 14 of a current sense resistor 16, a power field effect transistor (FET) 18, and three LEDs 20, 22, and 24. Transistor 18 has a drain 26, a gate 28, and a source 30. Transistor 18 also presents a gate to drain capacitance Cgd 32 and a gate to source capacitance Cgs 34, both indicated in dotted lines.\nIntegrated circuit controller 36 includes gate driver circuit 38, an Idrive output pin 40, and a PWM input pin 42. The Idrive output pin 40 is connected to the gate 30 of transistor 18 by external lead 41. Gate driver 38 includes switching circuit 44 connected between a power source and a ground and operates under control of a PWM signal received at pin 42 to turn on transistor 18 by sourcing charge to the gate 30 of transistor 18 and to turn off transistor 18 by removing charge from the gate 30. Turning on the transistor 18 conducts a current Iled from the power lead 12 through the resistor 16 and the LEDs 20, 22, and 24 to ground lead 14. Turning off the transistor 18 blocks current from the power lead 12 through the resistor 16 and the LEDs 20, 22, and 24 to ground lead 14.\nIn some applications, such as an automotive application, the resistor 16, power transistor 18, and LEDs 20, 22, and 24 are separate from the integrated circuit controller 36. In these applications, a problem exists in turning the current Iled through transistor 18 on and off in equal periods of time after state changes in the PWM signals. In large power FET transistors, the capacitances Cgd and Cgs are usually very large.\nThe current in the power FET transistors depends on the overdrive voltage VOV, which is the voltage difference between VGS and transistor GATE-SOURCE threshold voltage Vth. VOV_10% is the overdrive voltage when LED current reaches 10% of full current, VOV_90% is the overdrive voltage when LED current reaches 90% of full current, VOV_FULL is the overdrive voltage when LED has full current. The turn-on delay time Ton_delay depends on the time need to charge the capacitance Cgs from 0V to Vth+VOV_10%. The turn-off delay time Toff_delay depends on the time need to discharge the capacitance Cgs from VOV_FULL to VOV_90%. Typically the voltage difference between VOV_FULL−VOV_90% is much less than the difference between Vth+VOV_10% and 0V. This means that the turn-on delay time Ton_delay and the Toff_delay time of the transistor 18 will be different even with equal Idrive gate currents.\nIn FIG. 2, PWM signal turns on at time 46 and turns off at time 48. An Idrive current turns on at time 46 and turns off at time 48. Due to the capacitance Cgs 34, the voltage on the gate Vgs increases slowly to turn on the transistor 18 at time 50, effecting a Ton_delay between times 46 and 50. At time 48, the PWM signal turns off and the Idrive signal removes charge from gate 18, to turn off the transistor 18 at time 52, effecting a Toff_delay between times 48 and 52.\nA constant current Idrive is used to drive the transistor 18 in order to control the LED current slew rate during a rising and falling phase. For electromagnetic compatibility (EMC) considerations, a low slew rate (1˜10 mA/us) of the LED current is preferred. This means that the drive current cannot be very large. This results in the Ton_delay being long due to the limited drive current and a large Cgs of transistor 18 of about several hundred microseconds. The Toff_delay at time is very small, about several microseconds. A big gap between Ton_delay and Toff_delay exists, which causes Iled current duty cycle loss."}
{"text": "The present invention relates to vehicular traffic, particularly on long distance high speed highways, monitoring of the traffic, providing assistance to drivers in the traffic based upon the traffic monitoring and communication with specific vehicles in the traffic. The communications may originate from a vehicle, for example, identifying the vehicle and its location, may be sent to the vehicle to provide driving assistance, or may be sent to and received from the vehicle, for example, as in telephone communications. Further, the system provides for prioritizing travel on a multiple lane highway and for adjusting tolls charged for the use of the highway.\nCommunication with vehicles on high speed, long distance highways, monitoring traffic on the highways, and monitoring the positions and speeds of specific vehicles on the highways present substantially difficulties. In conventional mobile communications systems, for example, mobile telephones, fixed antennas are installed in the vicinity of highways. Usually, these antennas are elevated, for example, located on the tops of towers or buildings, in order to provide a large area of communication with vehicles. Each fixed antenna at least partially defines a cell and in typical cellular telephone communication, communication shifts from antenna-to-antenna, as a mobile transmitter moves between cells, usually without the notice of the persons, mobile or fixed, who are communicating.\nThe relatively widely spaced fixed antennas for cellular communication along highways have limitations. For example, each cell has a limited bandwidth from which channels for communication can be assigned. Thus, if too many telephone calls are attempted within a single cell at the same time, all channels may be placed in use so that some potential callers will not be assigned channels and will be unable to establish communication.\nIf traffic on a highway is to be monitored, and particularly if speeds and positions of individual vehicles are to be determined, simultaneous communication with each of the vehicles on the highway is required. Each vehicle requires a channel for communication. Absent a complicated multiplexing scheme, the bandwidth needed for communication within a typical mobile telephone cell between all of the vehicles traveling on a high speed long distance highway and a fixed antenna readily exceeds the available bandwidth. Therefore, such traffic monitoring is not even theoretically feasible. The bandwidth problem cannot be solved by increasing the available bandwidth because of the number of channels that would be required and limited electromagnetic spectrum availability.\nIt is an object of this invention to solve the problem imposed by the limited bandwidth available for communication with vehicles, particularly vehicles on a multiple lane high speed long distance highway, so that communication can occur with a large number of vehicles without the necessity of increased bandwidth of the communications.\nAccording to a first aspect, a traffic monitoring system includes lane terminals for detecting passage of a vehicle, a communication antenna, a terminal transceiver for communicating with a passing vehicle through the communication antenna, and a network backbone linking the lane terminals to a data processor for compiling information on passing vehicles sensed.\nIn a preferred arrangement, a traffic monitoring system for a highway includes first and second adjacent lanes for travel in the same direction, including a first line of the lane terminals located along an outside edge of the first lane, a second line of the lane terminals located between the first and second lanes, and a third line of the lane terminals located along an outside edge of the second lane, each of the first, second, and third lines of the lane terminals including respective network backbones connected the respective first, second, and third lines of the lane terminals.\nIt is particularly preferable that the system include at least one transverse link interconnecting the first, second, and third network backbones and a principal network backbone connected to the transverse link and providing an interconnection between the first, second, and third lines network backbones and the data processor.\nThe traffic monitoring system most preferable includes a traffic data base connected to the data processor through the principal network backbone for storing traffic information including passing vehicles detected by the sensor for processing by the data processor.\nThe traffic monitoring system provides for cellular communication with moving vehicles wherein groups of lane terminals define communication cells for communication with vehicles traveling on the highway and a cell management data base is connected to the data processor for identifying positions of specific vehicles on the highway with respect to the communication cells.\nFor increased utility, the traffic monitoring system may include a toll server connected to the principal network backbone and receiving information from the lane terminals for determining a toll of a vehicle traveling on the highway based upon the lane traveled by the vehicle.\nFor greatest utility, the traffic monitoring system includes mobile transceivers mounted on respective vehicles for sending signals to the lane terminals identifying the respective vehicle on which a transceiver is mounted.\nSimpler systems may omit communication antennas in the lane terminals or vehicle sensors in the lane terminals."}
{"text": "The disposal of used tires is increasingly becoming a difficult problem to solve due to increasing number of such tires which must be disposed of each year. In addition, the conventional means of burning tires is being precluded for environmental reasons. Therefore, there have been a number of devices disclosed in the prior art for reducing the disposal problems of tires.\nFor example, in U.S. Pat. No. 3,658,267 (Burwell), an apparatus for desintegrating tires and the like is disclosed. The apparatus includes a means for storing a plurality of tires therein and a means for feeding the tires one at a time to a holding means which rotates. A grinding and pulverizing means is disposed adjacent the holding means to reduce the tread of the tire to a pulverized form which can be removed and subsequently used. The remaining tire bead is then cut in two and a swing finger means is adapted to throw the tire bead down a chute.\nA tire detreading apparatus is disclosed in U.S. Pat. No. 1,745,021 (Liebau). This device includes a manually controlled device which is adjustable to receive tires of various sizes and to enable a quick and easy mounting of a tire thereon. Adjustment for tires of various sizes is provided by radially movable fingers which grip the rim of the tire and which are actuated by a disc including a series of spiral grooves. The cutter for the bead of the tire is also movable by manual operation into and out of engagement with the tire.\nOther U.S. patents which disclosed tire detreating apparatus or the like includes the following: U.S. Pat. No. 2,065,448 (George); U.S. Pat. No. 2,283,005 (Godfrey, Jr); U.S. Pat. No. 3,731,884 (Tupper et al); and U.S. Pat. No. 3,604,084 (Krieger)."}
{"text": "1. Field of the Invention\nThe present invention relates to the structure of three-dimensional image data, a method of recording three-dimensional image data, and a method of displaying and reproducing three-dimensional image data, and in particular, to the structure of three-dimensional image data, a method of recording three-dimensional image data, and a method of displaying and reproducing three-dimensional image data, wherein three-dimensional display images suitable for compression are recorded.\n2. Description of the Related Art\nVarious schemes are known for three-dimensional image display apparatuses that are able to display motion pictures three-dimensionally, that is, three-dimensional displays. In recent years, in particular, three-dimensional image display apparatuses of a flat panel type have been highly demanded which are based on a scheme not requiring any dedicated type of glass or the like. For flat panel type display apparatuses such as direct-vision or projection type liquid crystal display apparatuses and plasma display apparatuses whose pixels have fixed position in display surfaces, three-dimensional image display apparatuses can be easily implemented on the basis of the following scheme: a parallax barrier is installed immediately in front of the display panel to controllably direct rays from the display panel toward a viewer.\nLight rays are controlled so that even when observing the same position on the parallax barrier, the viewer views different images depending on the angle. Specifically, a slit array or lenticular sheet, that is, a cylindrical lens array, is used to provide only lateral parallaxes, that is, horizontal parallaxes, or a pinhole array or a lens array is used to also provide vertical parallaxes. The scheme using the parallax barrier is further classified into a two-view (binocular) scheme, a multiview scheme, a super-multiview scheme that is a multiview scheme provided with super-multiview conditions, and integral photography (hereinafter simply referred to as IP). The basic principle of these schemes is substantially the same as that which was invented about 100 years ago and has been used for three-dimensional photography since then.\nWith either the IP scheme or multiview scheme, viewing distance is normally finite. Consequently, a display image is created so that a perspective image is actually viewed at that viewing distance. An IP scheme providing only horizontal parallaxes and no vertical parallaxes is disclosed in SID04 Digest 1438 (2004). This one-dimensional IP scheme generates a set of parallel rays if the horizontal pitch of the parallax barrier is set equal to an integral multiple (n times) of the horizontal pitch of the pixels (the one-dimensional IP scheme is hereinafter also referred to as parallel-ray one-dimensional IP). Accordingly, a parallax component image in which pixel columns constituting a set of parallel rays are accumulated is obtained by perspective projection for a given viewing distance in the vertical distance and by orthographic projection in the horizontal direction. A parallax synthetic image (an elemental image array) is created by dividing each parallax component image obtained by the perspective projection in the vertical distance and by the orthographic projection in the horizontal direction, into pixel columns and interleaving and synthesizing. When the parallax image synthetic image is displayed on a display surface and observed through the parallax barrier, the correct projection, that is, a three-dimensional image based on perspective projection in both horizontal and vertical directions, is obtained. Specific methods are disclosed in SID04 Digest 1438 (2004), Jpn. Pat. Appln. KOKAI Publication Nos. 2003-90738 and 2003-315356, and the like. The multiview scheme divides an image obtained by a simple perspective projection, into pixel columns and interleaving and synthesizing to obtain a three-dimensional image based on the correct projection.\nIt is difficult to implement an image pickup apparatus that varies the projection method or a projection center distance between the vertical direction and the horizontal direction. This is because such an apparatus requires a camera or lens of the same size as that of a subject for orthographic projection. Accordingly, a practical method of obtaining orthographic projection data by image pickup is to convert perspective projection image pickup data into orthographic projection data. A known practical method is a ray space method based on interpolation and using EPI (epipolar surface).\nThe parallel-ray one-dimensional IP scheme is advantageous in that an image obtained by this scheme is easier to see than one obtained by the binocular scheme. However, the parallel-ray one-dimensional IP scheme requires complex image formats for the projection method and the dividing and arranging method. The binocular and multiview schemes are three-dimensional image display schemes that display the simplest three-dimensional images, and thus use simple image formats. With these schemes, each parallax image is formed of the same number of vertical pixels and the same number of horizontal pixels. For example, two parallax component images for the binocular scheme or nine parallax component images for the nine view scheme may be divided into pixel columns, which may then be synthesized into a parallax synthetic image in an image format in which the image can be displayed on the display surface. However, the parallel-ray one-dimensional IP scheme requires a larger number of parallax component images than the multiview scheme, which achieves substantially the same resolution. With the parallel-ray one-dimensional IP scheme, the number of horizontal pixels (horizontal range used to display three-dimensional images) in each parallax component image varies depending on parallax direction, resulting in complex image formats.\nBoth the multiview scheme and parallel-ray one-dimensional IP scheme pose a problem if each piece of parallax information is divided into subpixels assigned to a parallax synthetic image. When a parallax synthetic image is irreversibly compressed by an encoding method such as JPEG or MPEG, parallax information may be mixed together, and decompression may degrade image quality. In particular, if the color arrangement of color filters or the like are mosaic arrangement in order to avoid a moire effect, color information may also be mixed together, and decompression may further degrade image quality. Reversible compression avoids the degradation but offers a compression rate considerably lower than that achieved by the irreversible compression. A method exists which allows each parallax component image to be irreversibly compressed and decompressed. This method can be used with the multiview scheme but is inefficient with the parallel-ray one-dimensional IP scheme because it involves a large number of parallax component images and a varying number of horizontal pixels.\nAs described above, the conventional three-dimensional image recording method based on the parallel-ray one-dimensional IP scheme is disadvantageous in that it achieves a considerably low compression rate and may involve degradation of image quality as a result of irreversible compression and decompression."}
{"text": "The present invention relates to an audio encoding apparatus.\nAn audio encoding method that enables fast-forward play when reproducing data encoded by an MPEG (Moving Picture Expert Group) Audio Layer 3 (hereinafter referred to as MP3) method has heretofore been known.\nFor example, Patent Document 1 (Japanese Published Unexamined Patent Application No. 2006-190362) discloses an audio encoding method that quickly acquires information indicating a current play position and shortens a transition time from a normal play to a fast-forward or fast-backward play, and thus can avoid that a listener feels a temporary delay or stop in reproduction.\nIn this audio encoding method, an auxiliary data appender appends auxiliary data of 32×5 bits which are all defaulted to “0” to audio encoded data. An LBA writer overwrites a first 32-bit part of the auxiliary data with a value of LBA which is given by an LB counter. Further, a jump to position LBA writer overwrites the remaining parts of 32×4 bits with LBA_f4, LBA_f8, LBA_b4, and LBA_b8 to jump to which are given by the LB counter."}
{"text": "Semiconductor devices and semiconductor device fabrication processes are well known. For example, metal oxide semiconductor field-effect transistors (MOSFETs) are commonly used in a variety of different applications. As another example, LDMOS devices are often used in applications, such as smart power technologies, where higher voltage rating is desired and device size is at a premium. Indeed, power LDMOS devices are becoming increasingly popular for power applications. Furthermore, some power LDMOS devices must be designed to operate in a “high-side configuration,” which is a configuration in which all the device terminals are level shifted with respect to the substrate potential. A device that may be operated in a high-side configuration is said to be “high-side capable.” High-side capable power LDMOS devices are designed to prevent a direct punch-through path from a body region of the power LDMOS device to an underlying, heavily doped substrate.\nExisting technology attempts to satisfy high breakdown voltage requirements by utilizing a power LDMOS device having a reduced surface field (RESURF) structure. A power LDMOS device having a RESURF structure includes a first semiconductor region (which serves as a drift/RESURF region) having one conductivity type, and a second semiconductor region having a different conductivity type. The second region depletes the RESURF region from two different directions (bottom and side), allowing for a more uniform electric field distribution and thus reducing the peak electric field in the drift region, thereby allowing a higher breakdown voltage for the power LDMOS device. The RESURF structure just described is referred to herein as a “single RESURF” structure.\nA “double RESURF” LDMOS structure, on the other hand, includes first and third semiconductor regions having one conductivity type, and a second semiconductor region having a different conductivity type. In the double RESURF structure, first and third semiconductor regions deplete the second semiconductor region, thus reducing the peak electric field in the second semiconductor region to a greater degree than is possible with a single RESURF structure. Transistor devices, including power LDMOS devices and bipolar transistors, having single or double RESURF structures, will be referred to herein as “RESURF transistors.”\nLower “on” resistance and higher breakdown voltage characteristics are desirable in most practical LDMOS applications. An existing double RESURF LDMOS design utilizes a heavily doped n-type buried layer (NBL) that facilitates the double RESURF characteristics of the device. Unfortunately, the NBL limits the achievable breakdown voltage due to its depth relative to the source/body region near the surface of the device. Consequently, to achieve higher breakdown voltage, the thickness of the semiconductor material (e.g., p-type epitaxial semiconductor material) between the NBL and the source/body region could be increased. This solution, however, may not be feasible for higher breakdown voltage requirements, such as 80 volts and higher."}
{"text": "Known grid plates for seeding surfaces have cells each having a floor which is provided with an opening that is usually circular in shape. These openings are used, on the one hand, for draining the cells and, on the other hand, to enable root growth into earth below the grid plate. The purpose of the cell floors is to support the grid plate with respect to the earth, which is especially important when vehicles travel over the grid plates.\nThe openings at the bottom of the known grid plates are, however, not sufficiently large to ensure good rooting of grass with the soil. This results in a hindrance to the growth of the grass, drying out during long dry periods and the possibility of tufts of grass being torn out of individual cells.\nIn addition, it is also known to provide openings on the cell walls, as a result of which root growth between adjacent cells is possible. Although this increases the hold of the tufts of grass in the individual cells, the other aforementioned disadvantages remain."}
{"text": "EP 0 591 800 B1 discloses a fuel cell battery, which features in each case two interconnected plates arranged in parallel to each other and forming an intermediate cavity between two electrolytic electrode units arranged in parallel to each other and each consisting of an anode, a Polymer Electrolyte Membrane (PEM) and a cathode. Between the plates and the electrodes of the two adjacent electrolytic electrode units are embodied gas or fluid spaces for introducing operating media such as hydrogen or oxygen to the electrolytic electrode units.\nOne way in which the cavity between the two plates arranged in parallel to each other can be used is to feed in or to extract the cooling medium. In this case the two plates form a cooling card for cooling the two electrolytic electrode units. Instead of being used for inserting and removing the cooling medium, the cavity can also be used as a pressurized chamber for setting up a pressure pad which exerts pressure on the adjacent components and thereby creates a contact pressure between these adjacent components, to guarantee a good current flow between these two components for example. On the other hand such a pressure pad can be used as an element for compensate for tolerances in the dimensions of the adjacent components. However, when the cooling card is used as a pressure pad the cooling of the fuel cell block is reduced.\nA battery known from DE 27 29 640 C3 consisting of a number of electrochemical cells, especially fuel elements, features electrolytic electrode units arranged in parallel to each other, between which contact elements are arranged for tapping off power. These contact elements feature a cavity for feeding in a pressurized medium and are thus embodied as pressure pads for pressing together the overall arrangement. This enables an even contact pressure to be created between adjacent components and thereby for example the current flow between these components improved. No cooling of the electrolytic electrode units is provided here.\nCooling cards are additionally generally known, from WO 02/50953A2 for example. A cooling card disclosed in this application is embodied as a composite printed circuit board with two plates which form a cavity between them. This cavity features inlet and outlet channels and has cooling water or also heating water flowing through it during the operation of the fuel cells."}
{"text": "This invention relates to a method and apparatus for monitoring and diagnosing peripheral vascular blood flow in people or animals. Part of the apparatus uses a novel heat exchanger.\nPeripheral vascular disease (\"PVD\") is defined as impaired arterial or venous blood circulation. More than one percent of the United States population has some form of PVD.\nPVD may take many different forms, but it usually affects the leg, especially the femoral and popliteal (thigh and knee) regions. A major cause of venous PVD is deep vein thrombosis. Arteriosclerosis, which occurs more frequently with increased age, is usually progressive is a major cause of PVD in the arteries. Approximately forty percent of PVD cases are a result of complications from diabetes.\nAs PVD becomes more advanced, blood flow can be so limited that tissues will die, which leads to gangrene. As a result amputation is required above the gangreneous tissue. Cramps, charlie-horse or calf or thigh discomfort during leg movement are symptoms of milder PVD. Pain and the other symptoms increase as PVD becomes more advanced.\nThere are a number of methods for diagnosing PVD, but no simple screening devices for evaluating PVD accurately are now available. A number of the diagnostic techniques involve invasive diagnosis in which veins or arteries are entered by needles or catheters. In arteriography, after a contrast dye is injected into the arteries, the arteries are viewed using X-rays. Arteriography provides a detailed image of the arterial system, and it is the accepted standard for evaluation of an arterial disease before operations. Arteriography is expensive, and there is an element of risk, discomfort and exposure to radiation. Venography is similar to arteriography except the procedure is applied to veins. Both arteriography and venography are in widespread use but they cannot be used as a simple screening test for PVD.\nThere are also certain non-invasive tests for PVD in which the veins and arteries are not entered and are not exposed to X-ray examination. Non-invasive techniques have substantially less risks to the patient but they require skilled technicians and are usually limited to hospital settings. The non-invasive procedures now in use include segmental plethysmography, doppler ultrasound pulse analysis, electrical impedance measuring devices and a less-invasive arteriography called digital subtraction angiography. Segmental plethysmography measures blood volume in a limb or changes in blood pressure at different positions on the extremity such as at the thigh, calf, ankle and foot. A pulse volume recorder records arterial wave forms at each location. PVD in an artery alters the wave forms. As arteriosclerosis progresses, however, plethysmography becomes progressively less accurate. Therefore, an additional wave form analysis is performed on data provided by a doppler ultrasound instrument. As is well known, moving bodies shift the frequency of sound reflected from them. Doppler measuring instruments measure the difference between the frequency of transmitted sound and that of the reflected sound, and the difference is proportional to the velocity of the moving bodies. Red blood cells are such moving bodies, and their movement have been measured using doppler measurements. Plethysmography combined with doppler ultrasound measurements have provided a moderately accurate assessment of arterial PVD with less expense and lower risk than invasive techniques. Arteriography is still needed if surgery is contemplated. Plethysmography and doppler methods also measure the relative blood volume in a limb at a given time, and they are useful in evaluating deep vein thrombosis.\nElectrical impedance measuring devices have also been used to diagnose PVD. These devices operate on the principle that a decrease in the electrical resistance of a limb is directly proportional to the increase in the blood volume. Because of a relatively high occurrence of false, positive results, clinical use of these devices is declining.\nAlthough digital subtraction angiography is not invasive in that it does not require catheterization of the arteries, it does use radiation and the injection of a radiopaque tracer dye. A computer \"subtracts\" the X-ray picture after the dye is in the arteries from an X-ray taken before the dye is injected. Although this system is accurate, there is still a need for a simplified, non-invasive PVD diagnostic technique that involves no radiation or injection.\nIt is most important that the extent of PVD be diagnosed accurately and timely. In order to prevent gangrene, if PVD cannot otherwise be treated, amputation is frequently required. Obviously, the drastic step of amputation should not be undertaken unless PVD has progressed sufficiently. Therefore, false positive tests, giving incorrect evidence requiring amputation, should be avoided. The failure of a test to yield positive results when truly indicated is also a problem in that PVD can be life threatening. Once PVD is detected, preventative treatment can be timely started to avoid or delay amputation.\nOnce PVD has progressed beyond the stage that it can be treated to the point that amputation is required, a surgeon must decide where on the leg to amputate. If the amputation is done too low, PVD prevents sufficient blood flow to heal the amputation wound. If the amputation is too high, larger and more complicated prostheses will be necessary. This is especially true in deciding whether to amputate above or below the knee. Therefore, it is important to be able to determine the extent of PVD and its location.\nWith the exception of simple blood pressure testing and some doppler applications, the systems and methods discussed are performed in a hospital. This substantially increases their expense and decreases their general availability. What is desirable, therefore, is a simple, reliable and accurate test to measure vascular circulation in a physician's office (or even in a patient's home) both for initial screening and subsequent evaluations of PVD.\nIt is known that changes in the skin temperature of an extremity may indicate the presence or absence of blood flow. Heat transfer from the blood in the subcutaneous capillary bed to the skin is a function of the increase or decrease in blood supply from the arteries and the rate with which the blood is emptied from the capillary system by the venous system. Although factors such as muscle tone, activity, posture, venimotor tone, disease, reflexes and drugs affect cardiovascular function and blood flow, these affects can be minimized by environmental controls. Skin thermodynamics is therefore related to cardiovascular function insofar as cardiovascular function affects peripheral and capillary blood flow. Obstructions or occlusions in the peripheral vascular system will have a direct, measurable affect on heat transfer at the skin's surface.\nThe present invention may also have application in diagnosing all diseases that affect skin temperature and that now are diagnosed using infrared scanning."}
{"text": "1. Field of the Invention\nThe invention generally relates to devices for supporting loads and more particularly to a load-leveling base for supporting an adjustably positioned load such as a refrigerator of the type frequently employed in trailers, campers, motor homes, recreational vehicles and the like.\n2. Description of the Prior Art\nAs can be appreciated by those familiar with the design, fabrication and use of absorption refrigerators, liquid coolant circulated therethrough percolates to the top and then gravitates to the bottom of the cooling coil. If the refrigerator is not maintained in a leveled disposition percolation may be interrupted whereupon the cooling unit may sustain permanent damage. Hence, it is highly desirable that refrigerators of the absorption type be maintained in a leveled disposition throughout their operation.\nOf course, maintenance of a leveled disposition for a refrigerator within a recreational vehicle, such as a trailer, camper, motor home or the like often is achieved by first leveling the refrigerator in relation to the vehicle and thereafter leveling the vehicle upon parking. Because of the vehicles inherent mobility, operators frequently are burdened with a relatively difficult task since leveling of the vehicle is required upon each parking in order to assure that the refrigerator contained therein oriented in a properly leveled disposition. Hence, the utility of recreational vehicles of the aforementioned type tends to be impaired because of the time required to achieve a leveling of the vehicle, following each use thereof.\nIt is, therefore, the general purpose of the instant invention to provide an improved load supporting base adapted to be installed in recreational vehicles and employed for rapidly leveling and supporting loads such as absorption-type refrigerators, which is simple in design, economic to fabricate and practical to use."}
{"text": "1. Field of the Invention\nThis invention generally relates to controlling use of content, or other items, through rights expressions associated with the content or other items, and more particularly to a system and method for profiling rights expressions.\n2. Description of Related Art\nOne of the most important issues concerning the widespread distribution of digital content, such as documents, music, movies, software, information, and the like, in forms usable by computing devices, via electronic means, and the Internet in particular, is the provision of the ability to enforce the intellectual property rights during the distribution and use of the digital content. Technologies for resolving this problem are referred to as Digital Rights Management (DRM) herein. However, there are a number of issues to be considered in effecting a DRM system, such as authentication, authorization, accounting, payment and financial clearing, rights specification, rights verification, rights enforcement, and document protection issues, to name but a few.\nFor example, in the world of printed documents and other physical content, a work created by an author is usually provided to a publisher, which formats and prints numerous copies of the work. The copies are then sent by a distributor to bookstores or other retail outlets, from which the copies are purchased by end users. While the low quality of copying and the high cost of distributing printed material have served as deterrents to unauthorized copying of most printed documents, it is much easier to copy, modify, and redistribute unprotected digital content with high quality. Therefore, there is a need for mechanisms to protect digital content.\nDifficulties associated with preventing, or even deterring, people from making unauthorized copies of electronic content within current general-purpose computing and communications systems, such as personal computers, workstations, and other devices connected over communications networks, such as local area networks (LANs), intranets, and the Internet, are widely recognized. Many attempts to provide hardware-based solutions to prevent unauthorized copying have proven to be unsuccessful. Moreover, the deployment of high bandwidth or broadband communications technologies and the development of what is presently known as the National Information Infrastructure (NII) is making it more convenient to distribute large documents electronically, including video files, such as full length motion pictures, and this makes it easier to proliferate unauthorized copying and distribution of digital content. Therefore, the need for further development of DRM technologies is becoming a high priority.\nAccordingly, commonly-assigned U.S. Pat. No. 5,634,012 discloses a DRM system for controlling the distribution of digital content, wherein devices of the DRM system can include a repository associated therewith. A predetermined set of usage transaction steps define a protocol used by the repositories for enforcing usage rights associated with the content. Usage rights persist with the content and the usage rights associated with the content comprise a digital work. The usage rights can permit various manners of use of the content, such as a right to view or print or display the content, a right to use the content only once, a right to distribute or redistribute the content, and the like. Such usage rights can be made contingent on payment or other conditions. However, there is still a further need for systems and methods for expressing and enforcing usage rights and/or conditions associated with content in a flexible and robust manner."}
{"text": "This invention relates to methods for processing noisy sound signals, especially for nonlinear noise reduction in voice signals, for nonlinear isolation of power and noise signals, and for using nonlinear time series analysis based on the concept of low-order deterministic chaos. The invention also concerns an apparatus for implementing the method and use thereof.\nNoise reduction in the recording, storage, transmission or reproduction of human speech is of considerable technical relevance. Noise can appear as pure measuring inaccuracy, e.g., in the form of the digital error in output of sound levels, as noise in the transmission channel, or as dynamic noise through coupling of the system observed with the outside world. Examples of noise reduction in human speech are known from telecommunications, from automatic speech recognition, or from the use of electronic hearing aids. The problem of noise reduction does not only appear with human speech, but also with other kinds of sound signals, and not only with stochastic noise, but also in all forms of extraneous noise superimposed on a sound signal. There is, therefore, interest in a signal processing method by which strongly aperiodic and non-stationary sound signals can be analyzed, manipulated or isolated in terms of power and noise components.\nA typical approach to noise reduction, i.e. to breaking down a signal into certain power and noise components, is based on signal filtering in the frequency band. In the simplest case, filtering is by bandpass filters, resulting in the following problem however. Stochastic noise is usually broadband (frequently so-called xe2x80x9cwhite noisexe2x80x9d). But if the power signal itself is strongly aperiodic and thus broadband, the frequency filter also destroys a power signal component, meaning inadequate results are obtained. If high-frequency noise is to be eliminated from human speech by a lowpass filter in voice transmission, for example, the voice signal will be distorted.\nAnother generally familiar approach to noise reduction consists of noise compensation in sound recordings. Here, for example, human speech superimposed with a noise level in a room is recorded by a first microphone, and a sound signal essentially representing the noise level by a second microphone. A compensation signal is derived from the measured signal of the second microphone that, when superimposed with the measured signal of the first microphone, compensates for the noise from the surrounding space. This technique is disadvantageous because of the relatively large equipment outlay (use of special microphones with a directional characteristic) and the restricted field of use, e.g., in speech recording.\nMethods are also known for nonlinear time series analysis based on the concept of low-order deterministic chaos. Complex, dynamic response plays an important role in virtually all areas of our daily surroundings, and in many fields of science and technology, e.g., when processes in medicine, economics, signal engineering or meteorology produce aperiodic signals that are difficult to predict and often also difficult to classify. Thus, time series analysis is a basic approach for learning as much as possible about the properties or the state of a system from observed data. Known methods of analysis for understanding aperiodic signals are described, for example, by H. Kantz et al. in xe2x80x9cNonlinear Time Series Analysisxe2x80x9d, Cambridge University Press, Cambridge 1997, and H.D.I. Abarbanel in xe2x80x9cAnalysis of Observed Chaotic dataxe2x80x9d, Springer, N.Y. 1996. These methods are based on the concept of deterministic chaos. Deterministic chaos means that, although a system state at a certain time uniquely defines the system state at any random later point in time, the system is nevertheless unpredictable for a longer time. This results from the fact that the current system state is detected with an unavoidable error, the effect of which increases exponentially depending on the equation of motion of the system, so that after a relatively short time a simulated model state no longer bears any similarity with the real state of the system.\nMethods of noise suppression were developed for time series of deterministic chaotic systems that make no separation in the frequency band but resort explicitly to the deterministic structure of the signal. Such methods are described, for example, by P. Grassberger et al. in xe2x80x9cCHAOSxe2x80x9d, vol. 3, 1993, p 127, by H. Kantz et al. (see above), and by E. J. Kostelich et al. in xe2x80x9cPhys. Rev. Exe2x80x9d, vol. 48, 1993, p 1752. The principle of noise suppression for deterministic systems is described below with reference to FIGS. 10a-c. \nFIGS. 10a-c show schematically the dependence of successive time series values for noise-free and noisy systems (exemplified by a one-dimensional relationship). The noise-free data of a deterministic system produce the picture shown in FIG. 10a. There is an exact (here one-dimensional) deterministic relationship between one value and the sequential value. The time delay vectors, details of which are explained further below, lie in a low-dimensional manifold in the embedding space. Upon introduction of noise, the deterministic relationship is replaced by an approximative relationship. The data are no longer on the low-dimensional manifold but close to it as shown in FIG. 10b. The distinction between power and noise is by dimensionality. Everything leading out of the manifold can be traced to the effect of the noise.\nConsequently, the noise suppression for deterministically chaotic signals is made in three steps. First the dimension m of the embedding space is estimated and the dimension Q of the manifold in which the non-noisy data would be. For the actual correction, the manifold is identified in the vicinity of every single point, and finally the observed point is projected to the manifold for noise reduction as shown in FIG. 10c. \nThe disadvantage of the illustrated noise suppression is its restriction to deterministic systems. In a non-deterministic system, i.e., in which there is no unique relationship between one state and a sequential state, the concept of identifying a smooth manifold, as shown in FIGS. 10a-c, is not applicable. Thus, for example, the signal amplitudes of speech signals form time series that are unpredictable and correspond to the time series of non-deterministic systems.\nThe applicability of conventional, nonlinear noise reduction to speech signals has been out of the question to date, especially for the following reasons. Human speech (but also other sound signals of natural or synthetic origin) is very much non-stationary as a rule. Speech is composed of a concatenation of phonemes. The phonemes are constantly alternating, so the sound volume range is changing all the time. Thus, sibilants contain primarily high frequencies and vowels low frequencies. So, to describe speech, equations of motion would be necessary that constantly change in time. But the existence of a uniform equation of motion is the requirement for the concept of noise suppression described with reference to FIGS. 10a-c. \nIt is accordingly an object of the invention to achieve an improved signal processing method for sound signals, especially for noisy speech signals, by which effective and fast isolation of the power and noise components of the observed sound signal can be performed with as little distortion as possible.\nIt is also an object of the invention to provide an apparatus for implementing a method of this kind.\nA first aspect of the invention consists, in particular, in recording non-stationary sound signals, composed of power and noise components, at such a fast sampling rate that signal profiles within the observed sound signal contain sufficient redundancy for the noise reduction. Phonemes consist of a sequence of virtually periodic repetitions (forming the redundancy). The terms periodic and virtually periodic repetition are set forth in detail below. In what follows, uniform use will be made of the term virtually periodic signal profile. The recorded time series of sound signals produce waveforms that repeat at least over certain segments of the sound signal and allow application, on restricted time intervals, of the above mentioned, familiar concept per se of nonlinear noise reduction.\nAccording to another aspect of the invention, virtually periodic signal profiles are detected within an observed sound signal and correlations are determined between the signal profiles so that correlated signal components can be allocated to a power component and uncorrelated signal components to a noise component of the sound signal.\nYet another aspect of the invention is the replacement of temporal correlations by geometric correlations in the time delay embedding space, expressed by neighborhoods in this space. Points in these neighborhoods yield the information necessary for nonlinear noise reduction of the point for which the neighborhood is constructed.\nAnother aspect of the invention provides an apparatus for processing sound signals comprising a sampling circuit for signal detection, a computing circuit for signal processing, and a unit for the output of time series devoid of noise."}
{"text": "1. Field of the Invention\nThe present invention relates to a zoom lens system and an image pickup apparatus including the same, which is suitable for use in a television camera for broadcasting, a video camera, a digital still camera, and a silver-halide camera.\n2. Description of the Related Art\nIn recent years, there have been demanded a zoom lens system having a large diameter ratio, a high zoom ratio, and high optical performance for image pickup apparatuses such as a television camera, a silver-halide camera, a digital camera, and a video camera.\nA positive lead and telephoto type four-unit zoom lens system in which four lens units are provided in total and one of the lens units located closest to an object side has a positive refractive power has been known as the zoom lens system having a high zoom ratio. For example, U.S. Pat. No. 7,304,805 and Japanese Patent Application Laid-Open No. 2005-084409 disclose a four-unit zoom lens system which includes a first lens unit having a positive refractive power, a second lens unit having a negative refractive power, a third lens unit having a negative refractive power, and a fourth lens unit having a positive refractive power.\nU.S. Pat. No. 7,304,805 and Japanese Patent Application Laid-Open No. 2005-084409 disclose a four-unit zoom lens system in which an optical material having extraordinary dispersion is used to excellently correct a chromatic aberration, to thereby obtain high optical performance.\nWhen the optical material having extraordinary dispersion is used for the four-unit zoom lens system as described above, a lateral chromatic aberration and an axial chromatic aberration are excellently and easily corrected. In order to obtain high optical performance, it is important to excellently correct a lateral chromatic aberration at a wide angle end and an axial chromatic aberration at a telephoto end.\nFor example, in order to correct the lateral chromatic aberration at the wide angle end, the optical material having extraordinary dispersion only needs to be used for a lens which is located in an image plane side of an aperture stop and has a positive power (refractive power). However, in a case of an optical system having a relatively small F-number, the lens has no small effect on an on-axial beam, and hence the axial chromatic aberration tends to be overcorrected (axial chromatic aberration is excessively corrected).\nWhen the optical material having extraordinary dispersion is used to correct the lateral chromatic aberration at the wide angle end, the axial chromatic aberration tends to deteriorate in the entire zoom range. In order to increase the optical performance of the four-unit zoom lens system as described above over the entire zoom range, it is important to suitably set a lens structure of the fourth lens unit which does not move for zooming."}
{"text": "In many electronic applications, it is necessary to connect signal and/or power lines, for example optical signal lines to electronic devices in general, to other optical signal lines or to printed circuit hoards via terminals. In particular, when electronic devices are intended to be used outdoor, specific requirements have to be fulfilled by the connector assembly, in particular concerning water resistance and/or disconnection protection issues."}
{"text": "1. Field of the Invention\nThis invention relates to an apparatus for sealing a rotatable tubular member in a stationary housing which surrounds the end of the tubular member, comprising a gripping device, which is secured to the outside peripheral surface of the rotatable tubular member and rotates with the tubular member, and a sealing element, which is radially movably disposed in the gripping device and rotates with the latter and is in sealing relation to the housing.\n2. Discussion of Prior Art\nRotatable tubular members, particularly in rotary kilns, must be sealed at their charging and/or discharge ends in order to prevent or minimize an escape of gases and dusts and an ingress of air. For this purpose the gap between the rotatable tubular member and the stationary housing, which surrounds the end of the tubular member, must be sealed. As a sliding sealing element is required to seal the gap and in most cases is exposed to elevated temperatures, dust and gases, the sliding sealing element is subjected to heavy wear. Additionally, the rotatable tubular member is axially displaced as a result of thermal expansion and owing to the nature of its bearings and also performs continual radial movements because the axis of the rotatable tubular member never coincides exactly with the center line of the housing but performs a gyratory motion about the center line of the housing. Particularly these continual radial movements of the rotatable tubular member impose high stresses on the sliding sealing elements.\nThe periodical \"Zement--Kalk--Gips\" No. 5, 1971, contains on pages 208 to 215 a summary of arrangements for sealing rotary kilns. Whereas some of the sealing arrangements disclosed there are structurally expensive, none of them can prevent a heavy wear particularly as a result of the axial movement of rotary kilns.\nGerman Pat. No. 17 50 881 discloses a sealing arrangement comprising a tubular section which surrounds the outside peripheral surface of the rotatable tubular member and has on the inside an annular flange, which overlaps a ring that is secured to the outside peripheral surface of the rotatable tubular member. On the other side the flange is forced against the ring by thrust members, which are connected to the ring by screws. The flange and the tubular section are radially displaceable and are connected to the rotatable tubular member by chains so that the flange and the tubular section rotate with the rotatable tubular member. A sealing ring composed of a plurality of segments is secured to the stationary housing and surrounds and supports the tubular section. Owing to the asymmetrical loading of the tubular section, the gyratory motion of the end of the tubular member and the frictional contact between the flange of the tubular section and the gripping device, the radial displacement results in an additional load. The sealing segments always subjected to contact pressure wear heavily, the pressure contact of the sealing segments and the gripping device are highly liable to be deranged, and the sealing arrangement is highly expensive and can be mounted only with difficulty."}
{"text": "In recent years, the cost of traditional heating fuels such as heating oil and natural gas has significantly increased while at the same time periodic shortages or unavailability of such fuels has been experienced. As a result, many less conventional methods of space heating have grown in popularity correspondingly. One of the presently most popular alternative fuels is wood in that it is relatively inexpensive, plentiful and an efficient source of heat when burned. To make better use of firewood as a heat source, fireplaces and/or wood-burning stoves have grown greatly in demand as amenities for homes or similar buildings.\nTo facilitate the use of wood as a heating fuel, various types of log storage racks have been developed and marketed to provide a ready means of storage for a reasonable quantity of firewood. Such storage racks typically comprise little more than an open framework of metal tubing or the like adapted for stacking of firewood thereon. Many people consider firewood to be unsightly and out of place within one's home, even when stacked on a conventional storage rack and, for this reason, such storage racks have not met with widespread acceptance. Instead, many people will store firewood only outdoors necessitating constant trips to and from the outdoor storage area to continuously fuel a wood fire, which as will be understood is quite annoying and detracts from the desirability of wood fires as a heating medium.\nIn contrast, the present invention provides a log storage and dispensing apparatus of an aesthetically pleasing appearance in the nature of a furniture-type cabinet by which a quantity of firewood logs can be stored in a concealed manner indoors and the logs may be dispensed one-by-one as needed for fueling a fire."}
{"text": "1. Field of the Invention\nThis invention relates to fluid bed dryers and more particularly concerns a process and apparatus for injecting particles of material into a fluid bed dryer.\n2. Description of the Prior Art\nFluid bed dryers are used for drying particulate material by passing streams of fluidizing gases through a bed of material in a so-called fluidizing chamber. These chambers are usually separate removable units, frequently cylindrical in form and have perforate or screen-like bottom portions which have holes sufficiently small to retain the particulate material while permitting air or other fluid to flow therethrough. Fluid bed dryer operation involves passing sufficient air or other fluid through the chamber to suspend solid particles in a stream of air. If each particle is bathed with a stream of fluid such as air to produce the so-called \"fluid bed\", uniform drying throughout the entire mass is achieved and critical product qualities can be preserved.\nProduct loading operations of prior art fluid bed dryers are typically batch operations which involve the removal of the fluidizing chamber from the dryer, filling the chamber with particulate material, and replacing the chamber into the dryer. For some types of product, e.g., cheese and yeast, the individual particles readily adhere to one another and tend to agglomerate in the fluidizing chamber before the chamber can be returned to the dryer and the dryer activated. Problems of localized heavy air streams being channelled through the fluid bed chamber, product matting and non-uniform drying are increased by product agglomeration.\nOf particular interest to agglomeration problems is an improved fluid bed dryer as disclosed in U.S. Pat. No. 3,849,900. The improved fluid bed dryer of this patent is characterized by the use of a rotating apertured disc positioned below and out of contact with the product to be dried. The disc delivers streams of air moving in rotary fashion throughout the fluidizing chamber. This improved drying apparatus provides a mixing or agitation action in a fluidized bed of particulate material contained within the fluidizing chamber to free or break up particles which have agglomerated. The apparatus of the above-identified patent also eliminates localized heavy streams of air which tend to develop in the fluidizing chambers of prior art dryers and avoids mechanical agitation and contact with the particles which is undesirable for products that tend to crush, smear or agglomerate.\nAlthough the rotating disc dryer is a notable improvement over prior art fluid bed dryers, it only partially overcomes the problem of initial agglomeration which occurs when certain products are loaded into the product containing chamber of the fluid bed dryer in the conventional batch loading step. The initial agglomeration which occurs due to the batch loading of cheese, yeast or other soft particles frequently is not entirely broken up by the agitation in the ordinary operation of the rotating disc dryer. Thus, initial agglomeration results in nonuniformly dried products and a lack of uniformity in the particle size of the product due to the presence of agglomerates even in this improved prior art dryer."}
{"text": "The present invention relates to water-hydrolysable materials and, more particularly, to treatment fluids and associated methods relating to the hydrolysis of water-hydrolysable materials.\nWater-hydrolysable materials are commonly employed in subterranean operations. For instance, water-hydrolysable materials may be used in subterranean operations as fluid loss control particles, diverting agents, filter cake components, drilling fluid additives, cement additives, and the like. In some instances, the water-hydrolysable material may be in a mechanical form (e.g., plugs, sleeves, and the like). In another instance, the water-hydrolysable material may be capable of releasing a desirable degradation product, e.g., an acid, during its hydrolysis. The acid released by certain water-hydrolysable materials may be used to facilitate a reduction in viscosity of a fluid or to degrade a filter cake, as well as for numerous other functions in subterranean operations.\nInclusion of a water-hydrolysable material capable of releasing an acid in a gelled (and optionally crosslinked) treatment fluid may be used to facilitate a reduction in viscosity of such fluid. Generally, these water-hydrolysable materials likely will hydrolyze over time due to contact with water present in the fluid, thereby releasing an acid. Upon its hydrolysis, the acid will function, inter alia, to reduce the viscosity of the gelled (and optionally crosslinked) treatment fluid, for example, by breaking the crosslinks in the treatment fluid, reducing the pH of the treatment fluid sufficiently to reverse the crosslinks therein, and/or breaking down the backbone of the gelling agent present in the treatment fluid. Typically, the acid released by the water-hydrolysable materials may breakdown the gelling agents at temperatures above about 150° F.\nWater-hydrolysable materials capable of releasing an acid also may be used in the degradation of acid-soluble materials present in a subterranean formation, such as those present in or adjacent to filter cakes. Filter cakes commonly may be formed by a fluid (e.g., a drill-in and servicing fluid) on the face of a portion of a subterranean formation, inter alia, to minimize damage to the permeability thereof. The filter cake often comprises an acid-soluble component (e.g., a calcium carbonate bridging agent) and a polymeric component (e.g., starch and xanthan). Before desirable fluids, such as hydrocarbons, may be produced, the filter cake generally is removed. To facilitate the degradation of the acid-soluble component, a water-hydrolysable material capable of releasing an acid may be utilized. Filter cakes also may be removed using an acid where the filter cake does not contain an acid-soluble component, for example, by degrading the underlying carbonate adjacent, if the filter cake is present in a carbonate formation.\nIn one instance, the filter cake may be contacted by a treatment fluid that comprises the water-hydrolysable material. The resultant acid should interact with the acid-soluble component of the filter cake and/or the underlying carbonate adjacent to the filter cake in such a way as to facilitate their degradation. In another instance, the water-hydrolysable material capable of releasing an acid may be included in the fluid (such as the drill-in and servicing fluid) that forms the filter cake, such that the filter cake further contains the water-hydrolysable material. Subsequent contact of the filter cake with an aqueous fluid hydrolyzes the water-hydrolysable material thereby releasing an acid that acts to degrade the acid soluble component of the filter cake. Among other components, the aqueous fluid may contain oxidizing agents or enzymes suitable to facilitate the degradation of the polymeric component of the filter cake.\nUse of water-hydrolysable materials capable of releasing an acid may be problematic, for example, if the water-hydrolysable material hydrolyzes too slowly or too quickly. For example, where used to facilitate a reduction in viscosity of a treatment fluid, the treatment fluid may need to have a desired viscosity for a requisite duration to ensure a desirable well treatment. In some instances, as the temperature in the well bore increases, the hydrolysis rate of the water-hydrolysable material increases, which may lead to an untimely or undesired reduction in viscosity of the treatment fluid. One method used to reduce the hydrolysis rate of the water-hydrolysable material may be to encapsulate it in a slowly soluble coating that can delay the hydrolysis of the water-hydrolysable material and thus delay release of the resulting acid. However, encapsulation of the water-hydrolysable material may add undesired expense and complexity. Further, where the water-hydrolysable material has a relatively small particle size, e.g., less than about 200 microns, encapsulation may not be practicable. Also, while it is possible to “tune” the water-hydrolysable material through various methodologies (e.g., initial choice of material, choice of plasticizers, molecular weight of the material, etc.), these methods may not be sufficient to extend or decrease the degradation time appropriately and/or may not be economical."}
{"text": "It is known that abbreviations and acronyms are frequently used in the mobile telephony field. A glossary of acronyms/terms which are used throughout the present specification is set forth below:\nAPNAccess Point NameCSCircuit SwitchingGGSNGateway GPRS Support NodeGPRSGeneral Packet Radio ServiceGTPGPRS Tunneling ProtocolHLRHome Location RegisterIMSIP Multimedia SubsystemIMSIInternational Mobile Subscriber IdentityMSISDNMobile Subscriber Integrated Services Digital NumberMTMobile TerminatedP-CSCFProxy Call Session Control FunctionPDPPacket Data ProtocolPSPacket SwitchingRADIUSRemote Authentication Dial In User ServiceS-CSCFServing Call Session Control FunctionSGSNServing GPRS Support NodeU/SIMSubscriber Identity ModuleSIPSession Initiation ProtocolUMTSUniversal Mobile Telecommunications System\nAlways-on is an old and well known problem in the PS environment. Theoretically, any user who executes a PDP Context Request (makes a request for a new context) and receives a positive response, is using IP resources in an always-on manner, i.e. always present and assigned, even if the user does not generate any IP data traffic. The type of resources used has more of a static character (IP address, memory load in the nodes, etc.) than a dynamic character (radio channels, etc.), however even so it is a question of scarce resources. In fact current GGSNs are configured to close any PDP Context without activity for the purpose of saving the maximum number of resources as possible. The element defining which inactive PDP must be closed is the APN with which the access is carried out.\nWhen the IMS system is adopted on a large scale, it is possible that a paradoxical situation occurs. The use of the IMS for real time services implies an assumption of total reachability of the user at any time and in any place, with a view to maintaining the same level of service which the user currently has with CS. That means that the PDP Context for SIP signaling used by the user must always be available (from the IMS point of view, closing the signaling context means deregistering the user de facto), even if the inactivity of the user is extended for a long time. Therefore the IMS requirements will only worsen the problem of consuming IP resources without actual associated traffic.\nBefore explaining the proposed solution to this problem by means of the present invention, what would occur with the current situation is described below. It is assumed that there is an IMS user already registered in the system and with active signaling context, therefore he or she is capable of initiating and receiving calls. If the GGSN decides to close that signaling PDP Context because the user has not carried out any additional interaction with the network, the user would pass to a status similar to that of deregistered, being unable to initiate or receive any call until requesting a new context and registering again in the system. This implies additional signaling and furthermore the expected delay upon setting up any type of call increases."}
{"text": "The present invention relates to a method for treating and disposing of materials contaminated with blood or blood products and devices and materials for use in the process.\nMany diseases are transmitted between humans by contact of open wounds or sores with contaminated blood. Other diseases can be transmitted by insects, such as flies, fleas or mosquitoes, which feed on contaminated blood or materials which contain blood. These insects then contact the open wounds or sores on humans, bite or sting individuals or animals, or, as for example in the case of fleas on rodents, travel into living quarters of humans where the fleas can then transfer to humans, thus transmitting viruses, germs, bacteria, parasites or other disease vectors to the individual. Diseases known to be transmitted by contact with blood or which can be carried by insects are hepatitis, HIV, AIDS, malaria, herpes, human papilloma virus (HPV), and Jacobs Kreutzfelds disease. Many other diseases of unknown etiology are believed to be transmitted by insect bites or contact with infected blood.\nOver the last several years considerable efforts have been taken to eliminate or significantly reduce the possibility of transmission of blood bourn diseases by casual contact in the hospital or medical facility environment. Additionally, strict testing and controls have been placed on blood and blood products which are transfused during medical procedures. However, no attempts have been made to eliminate the possibility of the transmission of blood bourn diseases or the contamination of ground water supplies which can be caused by the casual or negligent disposal of menstrual blood. On a monthly basis, millions of women dispose of blood soaked tampons or pads in garbage receptacles or flush them down toilets. These blood soaked materials then end up in land fills or sewer systems where insects and other vermin can contact these blood soaked products. Unfortunately, if the woman's menstrual blood contains germs, bacteria, viruses, or blood parasites, these microorganisms can survive the disposal conditions and can then be transmitted to other humans. This is of particular concern in the case of HIV, hepatitis A, B, or C and HPV which are known to be transmitted in blood and which are believed to survive normal disposal conditions. Additionally, many diseases endemic to the less developed third world but virtually unknown in the developed countries of the world can be readily introduced into those developed countries by a single menstruating woman entering the country from those regions. This is a problem currently unrecognized or unaddressed by the world's public health authorities.\nThus, there is a need for a simple, safe and effective method and system to readily destroy microorganisms carried in menstrual blood and to allow safe disposal of menstrual blood."}
{"text": "As a communication standard of a serial bus for connecting between a personal computer and its peripheral device, and between Audio/Visual devices, there is the IEEE 1394 high performance serial bus standard (hereinafter, referred to as the IEEE 1394 standard).\nIn the IEEE 1394 standard, data transmission is specified at 100 Mbps (98.304 Mbps), 200 Mbps (196.608 Mbps), and 400 Mbps (393.216 Mbps), and a 1394 port having an upper transfer rate is specified to maintain compatibility with a lower transfer rate. Numeric values in the parentheses show actual transfer rates.\nThis allows a network to have the data transfer rates of 100 Mbps, 200 Mbps, and 400 Mbps.\nFurther, in the IEEE 1394 standard, as shown in FIG. 7, transmission data are converted into two signals, which are a DATA and a STROBE, which is supplementary to the DATA. The IEEE 1394 standard employs a transmission format of a DS-Link (Data/Strobe Link) coding method, which can generate a CLOCK from obtaining an exclusive OR of the two signals.\nAn arbitration signal at the physical layer is expressed with two sets of twisted-pair lines, namely TPA/TPA* and TPB/TPB*.\nTPA/TPA*, which is one of the sets of the twisted-pair lines sends a strobe signal (Strb_TX) and receives a data signal (Data_RX). On the other hand, another set of the twisted-pair lines, that is TPB/TPB*, sends a data signal (Data_TX) and receives a strobe signal (Strb_RX).\nA Strb_Tx signal, a Data_Tx signal, a Strb_Enable signal and a Data_Enable signal are used for generating arbitration signals (Arb_A_Rx, Arb_B_Rx) from a control code.\nHere, values and meanings of transmission arbitration signals are shown in FIG. 8, whereas values and meanings of receipt arbitration signals are shown in FIG. 9.\nFurthermore, in the IEEE 1394 standard, two kinds of connecting methods can be used, which are a daisy-chain and a node branching. In the daisy-chain method, a device including the 1394 port can be connected up to 16 nodes, and a distance between the nodes can be up to 4.5 m. Further, when the daisy-chain method is used together with the node branching, a network can be structured, which can have up to 63 nodes (physical node addresses) for connection, the maximum in the standard.\nFurthermore, in the IEEE 1394 standard, a cable can be plugged in or out while the device is operative, that is, while the power source is ON. When a node is added or deleted, a network is automatically reconstructed. At this time, a device connected to the node can be automatically recognized, and an ID or a location of the connected device is managed by an interface apparatus.\nIn data transmission between devices respectively including the interface apparatus that can respond to multiple transfer rates, speed negotiation, for example, is usually carried out as preparation of data transmission in a data transmission preparatory period. (Hereinafter, the devices are denoted as an own device and a counterpart device. The own device sends data to the counterpart device, and receives data from the counterpart device. The own device is provided with the interface apparatus discussed in the explanation, whereas the counterpart device is provided with the interface apparatus that corresponds to the discussed interface apparatus.)\nAs an example of the speed negotiation, there is a DAVIC (Digital Audio Visual Control) method. The speed negotiation method used in this method has a period for synchronizing bits and characters, a period for negotiating a communication speed, a period for confirming a termination of the speed negotiation, and a period for waiting for a constant time to retry the speed negotiation.\nIn the period for synchronizing bits and characters, first, a code having a large number of edges is sent in such a manner that the counterpart device that receives the code can easily synchronize bits, and then a PLL (Phase Locked Loop) in a receiving circuit synchronizes bits. After a predetermined period sufficient for a bit synchronization circuit to synchronize bits, a character synchronization circuit synchronizes characters. After synchronizing the characters, the period for negotiating the communication speed follows.\nIn the period for negotiating the communication speed, each of both devices between which the speed negotiation is performed respectively sends a code for indicating transmittable speeds of its own, and respectively receives a code for indicating a transmittable speed of the counterpart and thereby decides whether increase or maintain the transfer rate of its own at the present rate. After negotiation of the transfer rate is normally completed, the period for confirming the termination of the speed negotiation follows.\nOn the other hand, when the negotiation of the transfer rate is not completed within the predetermined constant time, the period for waiting for the predetermined constant time to retry the speed negotiation follows.\nIn the period for confirming the termination of the speed negotiation, the own device sends a code for indicating the termination of the speed negotiation and waits for receiving, from the counterpart device, a code for indicating the termination of the speed negotiation. When the own device receives the control code for indicating the termination of the speed negotiation, the speed negotiation is terminated, and then a period for data transmission follows.\nDuring the period for waiting for the constant time to retry the speed negotiation, a serial signal is not sent. At this point, because the receiving circuit of the counterpart device no longer receives a receipt signal, the receiving circuit detects that abnormality has occurred in the speed negotiation, and then the period for synchronizing bits follows.\nThe own device in a state for waiting for the constant time to retry the speed negotiation, after the predetermined constant time, moves to the period for synchronizing bits and resumes the speed negotiation.\nAs described above, by carrying out the speed negotiation, the own device and the counterpart device that respectively include the interface apparatus that can respond to multiple transfer rates can obtain a best possible transfer rate at which the both devices can transmit data in a channel having a low error rate.\nMoreover, the speed negotiation can be suspended and then resumed, even when the own device cannot normally receive the control code for indicating the transfer rate or the control code for indicating the termination of the speed negotiation, which is sent from the counterpart device, because a receiver of one device has abnormality or a channel has a high error rate.\nBy the way, a length of a metal cable in conformity to the IEEE 1394 is limited up to 4.5 m by reasons, such as transmission loss of a cable. At present, for long-distance transmission of a signal, an optical fiber, which has little transmission loss, is considered to be used in serial transmission, instead of the metal cable.\nWhen a signal of the IEEE 1394 standard is transmitted by using the optical fiber, the speed negotiation, for example, is carried out, as preparation of the data transmission in the data transmission preparatory period, thereby determining a best possible transfer rate at which the both devices can transmit data.\nIn the data transmission preparatory period, a receipt arbitration signal from an optical port to a state machine in conformity to the IEEE 1394 is no use, and electric power consumed for generating the arbitration signal during the period is wasted.\nFurthermore, during transition from the data transmission preparatory period to the data transmission period, the state machine in conformity to the IEEE 1394 may recognize an arbitration signal, which is being outputted into the state machine, as a valid arbitration signal depending on timing.\nIn such a case, a node which is not to be suspended may be unexpectedly suspended, causing a node state (a bus state) that is against user's intention, when the arbitration signal is an arbitration signal such as a SUSPEND or a DISABLE in conformity to IEEE 1394a-2000."}
{"text": "1. Field of the Invention\nThe present invention relates to a data information display method of a data broadcasting receiver and an apparatus thereof, and in particular to a data information display method of a data broadcasting receiver and an apparatus thereof which are capable of displaying data information corresponding to the present audition television picture as a simple data information picture without changing the size of the present audition television picture.\n2. Description of the Prior Art\nGenerally, in a digital data broadcasting, additional data information is transmitted from a transmitter to a receiver with an A/V (Audio/Visual) stream for a television by using digital signals.\nIn the digital data broadcasting, the transmitter compacts separately the video and audio, encodes it, multiplexes it, and transmits the data stream to the receiver. Then the receiver receives the data stream, corrects all errors, decodes it, and demultiplexes it into the original video and audio.\nIn order to watch the digital data broadcasting, an additional data processing S/W program for a digital broadcasting (hereinafter, it is referred to as a data browser) is required for the receiver in order to process the A/V stream and data and to support collectively various kinds of user support services, and the format and function of the data browser is similar to the format and function of an internet web browser.\nThe internet web browser is a S/W program based on the internet. It comprises a plurality of internet protocols, a HTML (Hyper Text Markup Language) document, a CSS (Cascading Style Sheets) Paser, a Java Script, a JAVA related engine, a formatting function for display, and its own program, and it supports most of e-mail protocols.\nIn addition, basic functions of the data browser used for the data information display apparatus of the data broadcasting receiver will now be described.\nThe data browser judges whether or not data is received related to the present audition broadcasting, detects the data when the data is received, informs the data receipt to a viewer by displaying an icon on a screen, and prepares a parsing, formatting, displaying of the received data information.\nAfter that, when the viewer wants to watch the data information receiving at the present and selects a “television and data broadcasting simulcast audition” item by using a remote-controller or an external input device, the data information picture is displayed. When the viewer does not watch the receiving data information, the present audition general television picture is displayed.\nFIG. 1 is a block diagram illustrating the conventional data information display apparatus of the data broadcasting receiver. It comprises a tuner 10 for tuning signals received through an antenna, an A/V decoder unit 20 for converting the signal tuned on the tuner 10 into an A/V stream, a screen 30 for displaying the A/V stream outputted from the A/V decoder unit 20, a data receiving unit 40 (listener) for searching whether there is the data information among the signals tuned on the tuner 10, and a PE (Presentation Engine) unit 50 for parsing and formatting the data information detected by the data receiving unit 40 and outputting the data information having the displayable format to the screen 30.\nHerein, the PE unit 50 comprises a user interface 57 for receiving an input signal of the viewer, a parsing unit 51 for receiving the user input signal though the user interface 57 and parsing the data information outputted from the data receiving unit 40, a formatting unit 53 for formatting the data information outputted from the parsing unit 51 in order to display, and a display unit 55 for processing the data information formatted on the formatting unit 53 so as to have the displayable format and transmitting it to the screen 30.\nHerein, the parsing unit 51 parses the data information outputted from the data receiving unit 40 by using the HTML (Hyper Text Markup Language) document, CSS (Cascading Style Sheet) Paser, Java Script, etc.\nThe operation of the data information display apparatus of the general data broadcasting receiver having the above-mentioned construction will now be described.\nFirst, the signal received through the antenna is tuned on the tuner 10, the tuned A/V signal is outputted to the A/V decoder unit 20, and the A/V decoder unit 20 converts the tuned signal into the A/V stream. The A/V stream is displayed on the screen 30.\nMeanwhile, the data receiving unit 40 detects whether there is data information among the signals tuned on the tuner 10, and stores it when the data information is detected. After that, when the viewer requests the data information output through the user interface 57 while the data receiving unit 40 receives the data information, the parsing unit 51 detects and parses the data information stored on the data receiving unit 40, and outputs the parsed data information. And, the formatting unit 53 formats the data information outputted from the parsing unit 51, and then the display unit 55 receives the data information formatted on the formatting unit 53, processes the formatted data so as to have the displayable format, and transmits it to the screen 30.\nFIG. 2 is a flow chart illustrating the general data information display method of the data broadcasting receiver. First, it is detected whether the data information is received among the signals tuned on the tuner 10 through the antenna (ST1). The general television picture is displayed when the data information is not detected (ST6). When the data information is detected, as depicted in FIG. 3A, the icon (for example, ‘i’) for indicating the presence of data information is displayed on the right upper portion of the screen 30 (ST2). At the same time, in order to output the received data information to the screen 30, the parsing, formatting, displaying are prepared (ST3). When the viewer selects the data information (ST4), in other words, when the viewer selects the “television and data information simulcast audition” icon through the remote-controller or external input device, as depicted in FIG. 3B, the data information is displayed on the screen 30 (ST5).\nMeanwhile, when the viewer does not select the “television and data broadcasting simulcast audition” icon, the present audition general television picture is continually displayed (ST6).\nFIGS. 3A and 3B illustrate the general data information display method of the data broadcasting receiver, where the construction of the picture can be changed according to the setting of the viewer.\nGenerally, because the switching of the picture is differently defined according to the received data information, when the data information picture is displayed, the size of the present audition general television picture decreases or only the data information picture is displayed. But, the viewer can check the size of the data information picture only after switching into the data information picture.\nIn addition, when the general data information display apparatus of the data broadcasting receiver generates the icon for indicating the data information receipt, before the viewer selects the “television and data broadcasting simulcast audition” item by using the remote-controller or external input device, the viewer can not know the content of the data information received at the present.\nIn addition, when the viewer selects the “television and data broadcasting simulcast audition” item by using the remote-controller or external input device, the present audition television picture is switched into the television picture and data information picture, wherein an interrupt for replacing the present audition television picture occurs because the size of the present audition television picture has to be adjusted again. Accordingly, the present audition television picture can not be displayed for a certain time.\nTherefore, when the present audition television picture is an important information, the viewer probably would not switch into the data information picture. Further, it is impossible to use the data information after a certain time. Accordingly the viewer can not check the data information finally."}
{"text": "Applications hosted in today's data centers suffer from internal fragmentation of resources, rigidity, and bandwidth constraints imposed by the architecture of the network connecting the data center's servers. The current conventional architecture statically maps web services to Ethernet VLANs, each constrained in size to a few hundred servers owing to control plane overheads. The IP routers used to span traffic across VLANs and the load balancers used to spray requests within a VLAN across servers are realized via expensive customized hardware and proprietary software. Expensive IP router ports thus become the bottleneck for any-rack-to-any-rack server connectivity, thus constraining the traffic for many data intensive applications (e.g., data mining, map/reduce computations, distributed file systems, blob stores). Further, the conventional architecture concentrates traffic in a few pieces of hardware that must be frequently upgraded and replaced to keep pace with demand—an approach that directly contradicts the prevailing philosophy in the rest of the data center, which is to scale out (adding more cheap components) rather than scale up (adding more power and complexity to a small number of expensive components). This concentration of traffic into a small number of network routers and switches also puts the network at risk of failures and outages, as failure of this small a number of components overcomes the redundancy built into the system and leads to an outage. Commodity switching hardware is now becoming available with very high port speeds at very low port cost, making this the right time to redesign the data center networking infrastructure.\nA data center is comprised of both server and networking components, where the distance between the components is typically less than a millisecond of speed of light propagation time (i.e., crossing a handful of switches at 1 Gbps speeds or greater). The server portion of the infrastructure is now far down the road of commoditization—high-end enterprise-class servers have been replaced by large numbers of low cost PCs. Innovation in distributed computing and systems management software have enabled the unreliability of individual servers to be masked by the aggregated reliability of the system as a whole. The running theme is “scaling out instead of scaling up,” driven by the economics of PC commoditization. Commodity parts are characterized by: wide availability, standardization, and drivers for increased capabilities (e.g., 1 Gbps, 10 Gbps technology already available, 100 Gbps technology now emerging).\nThe network portion of the data center infrastructure presents the next frontier for commoditization. The increase in the number of servers that need to be interconnected has stretched the limits of enterprise networking solutions so much that current architectures resemble a spectrum of patches and workarounds for protocols that were originally intended to work in enterprise networks orders of magnitude smaller.\nSome challenges and requirements of conventional data centers will now be explained with reference to FIG. 1, which shows a conventional architecture 100 for a data center, taken from a recommended source. See “Cisco systems: Data center: Load balancing data center services, 2004”, which is hereby incorporated by reference in its entirety. Multiple applications run inside the data center, but typically each application is hosted on its own set of (potentially virtual) server machines 102 with a single organization owning and controlling the activity inside the data center. Requests from the Internet 104 are typically sent to a publicly visible and routable Virtual IP address (VIP), and there are one or more VIPs associated with each application running in the data center.\nRequests arriving from the Internet are IP (layer 3) routed through border routers (BR) and access routers (AR) to a layer 2 domain based on the destination VIP address. The VIP is configured onto the two load balancers (LB) connected to the top switches (S), and complex mechanisms are used to ensure that if one load balancer fails, the other picks up the traffic. See “Virtual router redundancy protocol (VRRP)” by E. R. Hinden, which is hereby incorporated by reference in its entirety. For each VIP, the load balancers are configured with a list of Direct IP addresses (DIPs), which are the private and internal addresses of physical servers 102 in the racks below the load balancers. This list of DIPs defines the pool of servers that can handle requests to that VIP, and the load balancer spreads requests across the DIPs in the pool.\nAs the number of servers required in the data center grows, additional pairs of switches and associated racks are added to the layer 2 domain, as shown in the figure. Layer 2 subnets are constrained in size to a few hundred servers owing to the overheads of broadcast and control plane traffic, so VLANs are configured on the Layer 2 switches to divide up the domain into multiple layer 2 subnets, one sub-net per VLAN. When the layer 2 domains eventually hits limits associated with large Ethernet domains (e.g., VLAN exhaustion, broadcast/control traffic) at a size of about 4,000 servers, additional layer 2 domains are created and connected to other pairs of access routers.\nThe conventional approach has the following problems:\nFragmentation of resources: Popular load balancing techniques, such as destination NAT (or half-NAT) and direct server return, require that all DIPs in a VIP's pool be in the same layer 2 domain. See “Load Balancing Servers, Firewalls, and Caches” by C. Kopparapu, which is hereby incorporated by reference in its entirety. This constraint means that if an application grows and requires more servers, it cannot use available servers in other layer 2 domains—ultimately resulting in fragmentation and under-utilization of resources. Load balancing via Source NAT (or full-NAT) does allow servers to be spread across layer 2 domains, but then the servers never see the client IP, which is often unacceptable because servers need to log the client IP for regulatory compliance and data mining.\nPoor server to server connectivity: The hierarchical nature of the network means that for servers in different layer 2 domains to communicate, traffic must go through the layer 3 portion of the network. Since layer 3 ports are significantly more expensive then layer 2 ports, these links are typically oversubscribed (e.g., the capacity of the links between access routers and border routers is less than the sum of the output capacity of the servers connected to the access routers). The result is that the bandwidth available between servers in different parts of the data center can be quite limited. This creates a serious global optimization problem as all servers belonging to all applications must be placed with great care to ensure the sum of their traffic does not saturate any of the network links, and achieving this level of coordination between applications is difficult in practice. The lack of sufficient capacity between servers also fragments the pool of servers. For example, when an application running in the data center needs more servers to handle its workload, unused servers located elsewhere in the data center cannot be placed into service if there is insufficient capacity between them and the existing application servers.\nProprietary hardware that scales up, not out: The load balancers in the conventional architecture are used in pairs in a 1+1 resiliency configuration. When the load becomes too great for the load balancers, operators replace the existing load balancers with a new pair having more capacity, and it is impossible to add a single load balancer to obtain more capacity."}
{"text": "1. Field of the Invention\nThis invention relates generally to a fiber modified with carboxyl groups and, more specifically, to a cellulose or acrylonitrile fiber to which methacrylic acid or a hydroxyalkyl methacrylate is graft-copolymerized. The present invention is also directed to a process for producing such a modified fiber.\n2. The Prior Art\nOne known method of introducing carboxyl groups into a cellulose fiber includes reacting monochloroacetic acid with the cellulose fiber. Because of the introduction of carboxylic groups, the cellulose fiber thus obtained has improved hydrophilicity and is capable of absorbing basic, bad odor substances such as amines and ammonia. The known method has, however, a problem that it is difficult to homogeneously react the monochloroacetic acid with a solid, bulky mass of cellulose fibers because the reaction is exothermic.\nIt is also known to introduce ether linkages into a cellulose fiber by reaction with a vinyl compound according to the Michael reaction. With this method, however, it is not possible to effectively introduce carboxyl groups into the cellulose fiber when an unsaturated carboxylic acid such as acrylic acid or maleic acid is used as the vinyl compound. Thus, in order to introduce carboxylic groups into the cellulose fiber by utilizing the Michael reaction, it is necessary to first react the cellulose fiber with acrylonitrile and to hydrolyze the cyano groups into carboxyl groups.\nAn acrylonitrile fiber has properties similar to wool but has a defect that the moisture or damp absorbing power is poor. No effective method is however known to impart hydrophilicity to the acrylonitrile fiber.\nA method is known in which a vinyl monomer having an acidic group is graft-copolymerized to a polyester or polyamide fiber using am organic peroxide such as benzoyl peroxide."}
{"text": "The subject matter herein relates generally to a local luminaire area control system for outdoor lighting control.\nOn outdoor lighting, notably street lights and parking lot lights, photocontrol components and the corresponding mating receptacles are typically used to turn the lights on and off based upon the ambient light from the sun. Some light fixtures support dimming to variably control the light fixture based on the ambient light levels, time of day. There is a trend to provide programmable functions to the light fixtures based on sensors and programmable controls other than ambient light, such as, detected nearby pedestrian motion. To accommodate these functions, the lighting control receptacles provide both a light sensor and one or more environmental sensors in a sensor device on the light fixture. Different control systems require different mixes of functionality necessitating multiple product configurations having a different arrangement of sensors in the sensor device. When a different configuration is needed or desired a different sensor device is needed to replace the existing sensor device.\nA need remains for a sensor element that allows easy control and modification of sensor components in a sensor device."}
{"text": "This invention relates generally to electrodialysis and, more particularly to electrodialysis employing high brine concentrations.\nElectrodialysis is a membrane separation technology in which stacked pairs of selective cationic and anion selective membranes are typically used to segregate increasingly dilute salt streams from concentrated salt streams. Stacks of membrane pairs can be very large and can include 10 to 100 or more pairs of alternating membranes. At one end of the stack, electrochemical reactions are produced by a cathode in electrolyte solution. At the other end of the stack, another reaction is created by an anode in electrolyte solution. In the usual process, the electrolyte stream is separated from the dilute salt and the concentrated salt flows. The electrolyte solution is continuously applied to the electrodes. Furthermore, care is taken to collect the electrolyte streams from the anode and cathode reaction zones and remix the electrolyte prior to reapplication to the electrodes.\nAs will be appreciated by those skilled in the art, the general purpose of keeping solution flowing past the electrodes is to minimize the build-up of hydroxide ion at the cathode by virtue of acid destruction, and the build-up of excess acid at the anode by virtue of hydroxide destruction. Remixing the two electrolyte streams has the purpose of blend neutralizing the acid and base conditions that naturally occur.\nThe reactions at the anode and the cathode are integral to the process. It is the transfer of the electrons in the half-cell reactions that cause continued polarization that in-turn, causes the flux of salt ions across the selective membranes. Without initiation of the electrolysis reaction, electrodialysis cannot proceed.\nIn practice, there are several junctions within conventional electrodialysis processing that will or can impede efficient ion transfer. These include:\n1) Boundary diffusion resistance at all membrane surfaces;\n2) Resistance to current transmittal in all solutions;\n3) Relative difference in dilute and concentrate streams concentrations;\n4) Potential drop at each membrane; and\n5) Reaction initiation and transfer of ions at the electrodes.\nWhen electrodialysis is used in a conventional manner to treat light brine (e.g., brine that in general contains less than 1% salt and in some cases salt in a relative amount of as few as few hundred parts per million, i.e., high voltage and low amperage), one must be careful to avoid what is referred to as “limiting current” conditions. Limiting current is caused by applying too high a voltage such that instead of salt ion transfer, water at or near the membrane is caused to “split” into hydrogen ion and hydroxide ion. This can cause transfer of hydrogen ions and hydroxide ions across the respective membrane instead of the beneficial transfer of salt ions. Thus, amperage is generated by transfer of non-beneficial ions caused by the water splitting. This condition can be destructive to the membranes for a variety of reasons.\nHowever, when electrodialysis is used or applied at conditions to treat heavy brine or high salt concentrations (e.g., brine containing 1-8% salt down to 0.5% salt; i.e., low voltage and high amperage) such operation is very far from the “limiting current” conditions that pertain to conventional operation. Thus, such operation presents challenges and issues not addressed by conventional practice or operation.\nAn electrodialysis unit is typically composed of two fundamental components: electrodes cells consisting of an anode and a cathode; and a membrane stack of hydraulically isolated concentrate and diluate streams (the treatment cells) disposed therebetween. The electrode cells and the membrane stack are hydraulically isolated, but intimately associated electrochemically. Simplistically, the process can be viewed as an electric potential; pitted against a series of resistances; that generates an electric current at the electrodes. The electric current is translated into ion flux across the membrane stack. A partial list of the process dynamic and equilibrium resistances that can cause loss of potential (voltage drop) are summarized in the following Table 1, below.\nTABLE 1CathodeElectrode potential (Eo) for H+ reduction to H2CathodeDiffusion of H+ to the electrodeCathodeOvervoltage (unknown excess above Eo)/watersplitting at electrodeCathodeBack diffusion of anion from the electrode cellCathode/MembraneDiffusion of cation into electrode cellStackMembrane StackDiffusion of cation across the cation selectivediluate/concentratemembrane/Back diffusion of counter ionMembrane StackDiffusion of anion across the anion selectivediluate/concentratemembrane/Back diffusion of counter ionMembrane StackResistance of membraneDiluate/ConcentrateOsmotic pressure as a Nernst PotentialAnodeElectrode potential (Eo) for H2O/OH−oxidation to O2AnodeDiffusion of OH− to the electrodeAnodeOvervoltage (unknown excess above Eo)/watersplitting at electrodeAnodeBack diffusion of cation from the electrode cellAnode/Membrane StackDiffusion of cation into membrane stack\nElectrodialysis processing is conventionally operated in a salt concentration range of a few hundred ppm up to about 1% (10,000 ppm) total dissolved solids (TDS). Furthermore, conventional operation typically avoids cases where the divalent cation content of the water exceeds more than a few hundred milligrams per liter. Under such conditions, ion flux is usually limited by the “limiting current” at the membrane interfaces within the membrane stack. The voltage across the electrodes can be adjusted in order to avoid the limiting current condition. The limiting current is usually defined as a severe depletion of ions at the membrane surfaces such that the water boundary at the membrane “polarizes”. Effectively, the driving force to move ions (the stack voltage) exceeds the diffusion rate of ions to (and from) the membrane. Under polarized conditions, water “splits” into hydrogen ion and hydroxide ion. The ionized water is transported in lieu of the target ions. Thus, at this limiting current condition, power is wasted by transferring hydrogen ion instead of the target cation (example: sodium ion) and/or hydroxide ion instead of the target anion (example: chloride). The polarized water also can cause rapid degradation of the membrane materials. Therefore, under these conventional salt ranges, the majority of development in the science of ED technology is driven by the prediction and avoidance of the polarization of water at the membrane surfaces.\nIn such conventional processing, the chemistry of the electrolyte is relatively simple. For example, a typical electrolyte solution contains about 30 g/L Na2SO4 and has a conductivity of about 28-30 mS/cm. This is usually sufficient since the water to be treated, 1% TDS as NaCl or less, has a conductivity of less than 16 mS/cm.\nA steady flux of ions across the membrane stack balances the electrochemical reaction at the electrodes. Hydrogen ions are reduced to hydrogen gas at the cathode. The loss of H+ in the catholyte causes an ion imbalance that requires the transport of a cation (e.g., sodium ion) from the membrane stack into the cathode cell. Oxygen is produced at the anode at the expense of hydroxide ion. The resultant loss of hydroxide by such oxidation necessitates the transport of a cation (e.g. sodium ion) away from the anolyte into the membrane stack in order to appease the loss of hydroxide at the cathode. The result is that the catholyte becomes basic and sodium rich while the anolyte becomes acidic and sodium poor. In order to maintain pH and sodium balance between the catholyte and the anolyte, it is common (after suitable degassing) to re-mix catholyte and anolyte to make a common electrolyte for re-application to the electrodes.\nHowever, the application of electrodialysis to high brine conditions, i.e., 0.5% to 8% TDS, presents many and various challenges. For example, with the application of electrodialysis to high brine conditions the opportunities to create a “limiting current” condition at the membrane surfaces are typically severely limited as the ionic strength of the solution is so high.\nMoreover, in highly concentrated brines containing significant concentrations of multivalent cations, such as magnesium, calcium, barium, and iron, the cations can cross from the water (usually from the diluate cell located adjacent to the cathode cell) into the cathode cell. Precipitation of calcium and barium as hydroxides or sulfates, or the precipitation of iron as a hydroxide at the electrodes is problematic and can result in a rapid increase in resistance to ion flux at the surface of the electrode. In practice this is observed as a decrease in current (amps) for a given voltage. In conventional electrodialysis, the presence of such multivalent cation ions cause a need for frequent cleaning of the electrodes. Conventional cleaning usually involves alternating flushes of both the electrode cells and the treatment cells with a strong acid (e.g. 1% HCl) followed by strong base (e.g., 1% NaOH). Conventional washing typically requires significant down time. An additional disadvantage of conventional cleaning is the need for handling and storage of strong chemical cleaning agents.\nFurther, while pole reversal has previously been applied during conventional electrodialysis, such conventional practice has typically involved switching the polarity of the electrode for extended periods of time, up to several hours per switch. In order to accommodate such pole reversal, there must also be a hydraulic shift that switches the concentrate stream to become the diluate stream, and vice-versa.\nIn view of the above, there is a need and a demand for improved electrodialysis processing and, more particularly to improved electrodialysis processing when using or in conjunction with the employment of high brine concentrations."}
{"text": "Reclosable containers are convenient and easily loadable packages which are advantageously used when it is desired to repeatedly access a product within a container, or more commonly when a package with a positive locking closure is needed which may be easily loaded and sealed without the need for adhesives or sealing tools. These containers are commonly thermoformed of thin thermoplastic material and have a bottom portion or base and a top lid or cover. To keep the container closed after initial opening, it is necessary to provide some means of interlocking the cover and the base. A common type of interlocking structure employs projections and depressions on shoulders extending from the cover and base which engage with each other in a \"snap\" fit. The cover and the base are thus held together during normal usage but may be pulled apart by applying sufficient force.\nTo avoid unwanted opening of the container, it is advantageous to have a tight seal between the cover and the base. However, if the seal is too tight it will require excessive force to reopen the container. This difficulty in opening is compounded by the fact that with a standard snap fit closure the separating force must be applied by pulling the top cover and base in opposite directions. With containers that are too large or too small to be easily grasped, even greater difficulties in opening will be encountered.\nThe problem of producing a container with a tight reliable seal, yet that is easy to open, is long standing in the art and has evoked many different approaches to solution. What is needed is an easily manufactured package that may be securely closed yet easily opened."}
{"text": "1. Field of the Invention\nThe invention pertains to the field of tensioners. More particularly, the invention pertains to a self-energizing brake for a tensioner.\n2. Description of Related Art\nChain drive applications require a tensioning device to maintain tension to control the motion of the chain. Blade spring tensioners are often used to tension a chain drive, and conventional blade spring tensioners apply similar damping forces in both directions of travel. Increased damping is desired when the chain tensioner is being compressed by the chain drive, and then when the tensioner is extending to take up slack in the system.\nU.S. Pat. No. 5,935,032 discloses a wheeled tensioner for belts. As the tension from the belt presses on the wheeled tensioning pulley, a first pivot arm having an end with an arcuate surface, compresses a spring that moves a second pivot arm having an end with a movable shoe, and creates friction between the shoe convex arcuate friction surface and the arcuate surface of the first pivot arm.\nU.S. Pat. No. 6,699,148 discloses a self adjusting contact mechanism for a belt that includes a contact shoe with a low friction surface with a belt, and a mechanism for biasing the low friction surface of the shoe against the belt with a controlled force profile. The shoe is mounted to a lever arm through a first pivot joint. The first pivot joint maintains the shoe in alignment with and in parallel with the belt. The second pivot connects a lever onto a belt drive system.\nJapanese Publication No. 06193693 discloses a blade tensioner in which at the proximal end is a rocker fulcrum and a wedge type interspace formed in an interval between a cam surface and a cylindrical bore. A friction member and a spring member are built in each wedge type interspace. When the tension in the chain increases the friction member is thrust into a narrow part of the interspace and any rocking motion of the chain guide is prevented.\nJapanese Publication No. 2001-355688 discloses a first embodiment in which a friction material is inset into the distal end of the blade shoe and contacts the sliding surface of the bracket, increasing dampening. In another embodiment, a friction surface is inset into the sliding surface of the bracket. In an alternative embodiment, the friction surface is inset into the curled ends between the end of the blade shoe and the blade spring.\nJapanese Publication No. 08-336809 discloses an energized brake that stops rotation of a chainsaw blade when the motor is turned off. The brake device includes a brake band surrounding a brake drum with an end of the brake band operable by a hand guard, so that when activated, the brake drum is clamped by the brake band, stopping the chain. At the same time, a wire energizes a coil spring contacting a brake shoe, so that when a switch is turned off, the wire energizes the coil spring and thus the brake shoe, causing the shoe to engage the brake drum. When the switch is turned on, the brake shoe resists the coil spring and the shoe does not engage to brake drum.\nTherefore, there is a need in the art for a self-energizing brake that increases damping of the blade tensioner during compression of the blade shoe rather than during extension of the blade shoe."}
{"text": "A customer of a long distance telephone service provider may be charged for long distance service based on the features of each long distance call such as the time of day of the call, the length of the call, the destination of the call, etc.\nAlternatively, the customer may choose to take advantage of a calling plan offered by the service provider. For example, the customer may subscribe to a calling plan that provides him or her with a set period of time of long distance service during a particular month (or other time measure) at a specified cost. The customer pays the specified cost even if the customer does not use up all of his or her allowed time in the month. On the other hand, if the customer uses up more than the allowed time, the customer may be subject to extra or even premium charges. In some cases, the customer may be prevented from using more than the set period of time. If, for example, the customer is in the middle of a long distance call and the set period of time runs out, then some service providers may cut-off the call.\nAs yet another alternative for obtaining long distance telecommunication services, a customer may purchase a pre-paid calling card. Unlike the typical calling plan described above, the customer may make use of the pre-paid calling card from almost any telephone number. Also unlike the typical calling plan, the customer generally does not have to use up the time allowed by the pre-paid calling card within a month or some other limited amount of time. Nonetheless, there may be some expiration date of a pre-paid calling card. Like the calling plan, the pre-paid calling card provides the customer with a set period of time of service for a specified cost.\nAs a result of the limitations of calling plans and pre-paid cards, the customer may strive to keep his or her use of long distance service to just under the set period of time allowed by the calling plan or pre-paid calling card. If the customer uses too much time, there are financial consequences or even an interruption in service. But if the customer uses too little time, then the customer may not be getting the most value out of the calling plan or pre-paid calling card.\nThere are difficulties, however, in determining how much time of a calling plan or a pre-paid calling card has been used up or is available. The customer could keep track of the allowed time. Whenever the customer made a call, the customer could time the length of each call and accumulate the call lengths to determine how much of the allowed time has been used up. But having to time each call is inconvenient, inefficient, and distracting. Moreover, people other than the customer, like the customer's family, may be making calls. It would be difficult to keep track of everybody's calls.\nAnother way to keep track of the time available or used in a calling plan or pre-paid calling card is to contact the service provider. Obviously, the service provider keeps track of allowed time so as to cut-off a customer who has exceeded a limit or to assess extra charges. Typically, the service provider designates a particular telephone number to be used by a customer to find out the balance of time remaining or time used up. The customer calls the particular telephone number, and may be requested to enter his or her telephone number subject to the calling plan. In the case of a pre-paid calling card, the customer may have to enter the number assigned to the pre-paid calling card. In addition, the customer may have to provide identification such as a password or personal identity number (PIN). Thus, in addition to the 7 or 10 digits necessary to make the call to the service provider, the customer may have to enter another 7 or 10 digits to identify his or her telephone number or pre-paid calling card, and may have to enter additional digits to provide a password or PIN. Moreover, the customer has to repeat this process each time the customer desires to know how much time may be left in his or her calling plan or pre-paid calling card. This method of keeping track of the available time in a calling plan or a pre-paid calling card is cumbersome, wastes time, and is inconvenient.\nAccordingly, there is a need for providing information to a customer about his or her calling plan. In particular, there is a need for providing a customer with information about the time remaining or used-up in a calling plan. There also is a need to provide the customer with the calling plan information in a convenient fashion so the customer does not have to keep track of the allowed time himself or herself, does not have to make separate calls to a service provider, and does not have to provide a lot of information as a requisite to obtaining the calling plan information."}
{"text": "The present invention relates generally to seat restraint systems for vehicles and, more particularly, to a selectable load limiting seat restraint retractor for a seat restraint system in a vehicle.\nIt is known to provide a seat restraint system such as a seat belt in a vehicle to restrain an occupant in a seat of the vehicle. In some vehicles, the seat restraint system may be a lap restraint or belt, a shoulder restraint or belt, or both. Typically, the lap belt and shoulder belt are connected together at one end. The seat restraint system includes a latch plate at the connected end. The seat restraint system also includes a buckle connected at one end by webbing or the like to vehicle structure. The buckle receives the latch plate to be buckled together. When the buckle and latch plate are buckled together, the seat restraint system restrains movement of the occupant to help protect the occupant during a collision.\nIt is also known to provide a load limiting retractor for the seat restraint system. During a front impact of the vehicle, the load limiting retractor locks and rapid deceleration causes forward movement of the occupant and the occupant to load the seat restraint webbing. These load limiting retractors are used to payout or unspool seat restraint webbing under a constant load and allow the occupant to decelerate with the vehicle, thus limiting high seat belt restraint loads, thereby improving occupant performance. However, these load limiting retractors offer only one level of load limiting and are usually designed for mid sized occupants.\nTherefore, it is desirable to limit the amount of restraint load on an occupant during an impact for a seat restraint system in a vehicle. It is also desirable to select the amount of load limiting based on a size and mass of the occupant during an impact for a seat restraint system in a vehicle. It is further desirable to provide multiple levels of load limiting such as both a high level and a low level for various size occupants during an impact in a seat restraint system of a vehicle. Thus, there is a need in the art to provide a selectable load limiting seat restraint retractor that meets these desires.\nIt is, therefore, one object of the present invention to provide a new selectable load limiting seat restraint retractor for a seat restraint system in a vehicle.\nIt is another object of the present invention to provide a selectable load limiting seat restraint retractor for a seat restraint system in a vehicle that has multi-level load limiting.\nIt is yet another object of the present invention to provide a selectable load limiting seat restraint retractor for a seat restraint system in a vehicle having both a high level and a low level of load limiting.\nTo achieve the foregoing objects, the present invention is a selectable load limiting seat restraint retractor for a seat restraint system in a vehicle including a first housing for operative connection to vehicle structure. The selectable load limiting seat restraint retractor also includes a rotatable take-up spool connected to webbing of a seat restraint system for winding and unwinding the webbing. The selectable load limiting seat restraint retractor further includes a rotatable first torsion bar operatively connected to the first housing to provide a first level of load limiting on the webbing during an impact. The take-up spool is operatively connected to the first torsion bar. The selectable load limiting seat restraint retractor includes a rotatable second torsion bar operatively connected to the first housing and orientated parallel to the first torsion bar to provide a second level of load limiting on the webbing during an impact. The second level of load limiting is less than the first level of load limiting. The selectable load limiting seat restraint retractor further includes a selectable load limiting mechanism operatively cooperating with the first torsion bar and the second torsion bar to select between the first level of load limiting and the second level of load limiting to provide digressive load limiting on the webbing during an impact.\nOne advantage of the present invention is that a selectable load limiting seat restraint retractor is provided for a seat restraint system in a vehicle having a digressive load limit. Another advantage of the present invention is that the selectable load limiting seat restraint retractor achieves multi-level load limiting and provides the desired effect of digressive load limiting. Yet another advantage of the present invention is that the selectable load limiting seat restraint retractor offers a high peak load level early in an impact prior to airbag contact and digresses to a lower load level as an occupant is restrained by the airbag, thereby spreading the impact loading more evenly throughout the impact. Still another advantage of the present invention is that the selectable load limiting seat restraint retractor allows larger sized occupants to benefit by receiving a high load limit level throughout the impact and from the additional load carrying capability of the seat restraint system. A further advantage of the present invention is that the selectable load limiting seat restraint retractor allows smaller sized occupants to benefit from a lower load limit because these occupants are in most cases fully restrained by the airbag and little additional restraint from the webbing is necessary. Yet a further advantage of the present invention is that the selectable load limiting seat restraint retractor maintains critical packaging dimensions. Still a further advantage of the present invention is that the selectable load limiting seat restraint retractor has a high level of load limiting as a default with a low level of load limiting being selectable during an impact. Another advantage of the present invention is that the selectable load limiting seat restraint retractor has a load level response that starts high and ends lower, working in conjunction with the airbag to maximize performance for a mid sized occupant in a high severity impact.\nOther objects, features, and advantages of the present invention will be readily appreciated, as the same becomes better understood, after reading the subsequent description taken in conjunction with the accompanying drawings."}
{"text": "This invention relates to a ceramic membrane device which is suitable for a photomask and to a method for producing the ceramic membrane device.\nWhen processing an integrated circuit, a photomask is used for the purpose of forming a pattern of electronic circuits or masking a pattern on a wafer. A conventional ceramic membrane device is composed of a ring-shaped base plate made of silicon and a circular ceramic coating supported on the base plate.\nIn the conventional ceramic membrane device, the ceramic coating is formed on only one face of the base plate. Thus, the base plate is apt to be disadvantageously curved or bent because there is a difference in magnitude of thermal expansion between the silicon base plate and the ceramic coating. If the base plate is curved, the ceramic coating is also curved so that its flatness deteriorates. In case of a curved ceramic membrane on which a masking pattern is formed, the pattern formed on a wafer by means of an X-ray will be faulty."}
{"text": "This application is an application under 35 U.S.C. Section 371 of International Application Number PCT/FR991/01691 filed on Jul. 9, 1999.\nA subject matter of the present invention is a novel terephthalic polyester composition, its process of preparation by transesterification/condensation of poly(ethylene terephthalate) and of polyethylene glycol, and its use as soil-release agent in detergent formulations, in rinsing, softening or finishing formulations, for the washing, with or without pretreatment, the rinsing, the softening or the finishing of textiles, in particular polyester-based textiles.\nThe soil-release activity of ethylene terephthalate/polyethylene oxide terephthalate copolymers in the finishing of textiles, in particular polyester-based textiles, and the use of said copolymers as soil-release and/or antiredeposition agents in detergent formulations for the washing, with or without pretreatment, of textiles, in particular polyester-based textiles, are well known (U.S. Pat. Nos. 396,252; 4,116,885; 4,785,060).\nThese copolymers can, for example, derive from the transesterification/condensation of poly(ethylene terephthalate) and of polyethylene glycol (U.S. Pat. No. 4,785,060).\nThe Applicant Company has found a novel terephthalic polyester composition which can be obtained by transesterification/condensation of a poly(ethylene terephthalate) and of a polyethylene glycol and which in particular exhibits particularly good soil-release properties.\nAccording to the invention, it is a terephthalic polyester composition (TPC) comprising, as a mixture\nan ethylene terephthalate homooligomer (PET1) essentially comprising oxyethylene terephthalate (TE) repeat units of formula (I)\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94xe2x80x83xe2x80x83(I)\nxe2x80x83where A represents the 1,4-phenylene group,\nand a block terephthalic copolymer (PET2/TE-POE) comprising\nat least one polyethylene terephthalate block (PET2) composed of oxyethylene terephthalate (TE) repeat units of formula (I)\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94xe2x80x83xe2x80x83(I)\nxe2x80x83where A represents the 1,4-phenylene group,\nand at least one polyoxyethylene terephthalate block (TE-POE) of formula\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94(CH2xe2x80x94CH2xe2x80x94O)nxe2x80x94\nxe2x80x83the value of n being such that said block exhibits a number-average molecular mass of the order of 1500 to 4000, preferably of the order of 3000 to 4000,\nsaid composition being characterized in that:\nthe amount of (TE) units of the polyethylene terephthalate (PET1) does not represent more than 10%, preferably not more than 7%, of all the (TE) units present in the terephthalic polyester composition (TPC),\nthe amount by weight of all the (TE) units present in said (TPC) composition represents at least 11%, preferably from 11.5 to 17%, of the weight of said (TPC) composition,\nthe amount by weight of mono(oxyethyleneoxy) (OEO) residues of formula\n1/2Oxe2x80x94CH2xe2x80x94CH2xe2x80x94O1/2\nxe2x80x83represents at least 1.3%, preferably from 1.3 to 2.3%, of the weight of said terephthalic polyester composition (TPC),\nsaid (OEO) residues belonging to the oxyethylene aromatic diester (OAD) groups of formula\nxe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94Oxe2x80x94C (O)xe2x80x94Axe2x80x94\nxe2x80x83present in all the (PET1) and (PET2) blocks,\nand in that the weight-average molar mass of said block terephthalic copolymer (PET2/TE-POE) is at least 30000, preferably at least 35000, very particularly at least 40000.\nThe ends of the chains of (PET1) homooligomers are generally composed of xe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94OH units.\nThe ends of the chains of block terephthalic copolymer (PET2/TE-POE) are generally composed of\nxe2x80x94C(O)xe2x80x94Axe2x80x94C(O)xe2x80x94Oxe2x80x94CH2xe2x80x94CH2xe2x80x94OH\nunits and/or\nxe2x80x94(CH2xe2x80x94CH2xe2x80x94O)nxe2x88x921xe2x80x94CH2xe2x80x94CH2xe2x80x94OH\nunits.\nThe characteristics of the (TPC) composition of the invention can be determined by subjecting said composition to the following analyses:\nsize exclusion chromatography (SEC) using a size exclusion chromatography device with simultaneous detection by refractometry and ultraviolet, using N,N-dimethylacetamide (DMAc) comprising 10xe2x88x922 mol/l of LiBr as eluent at 100xc2x0 C.\nThis measurement (UV chromatography) makes it possible:\nto detect the presence of the A (terephthalic) aromatic nuclei,\nto determine the percentage of ethylene terephthalate homooligomer (PET1) in the terephthalic polyester composition (TPC),\nto determine the weight-average molar mass of the block terephthalic copolymer (PET2/TE-POE), expressed in polystyrene equivalents.\nproton nuclear magnetic resonance (1H NMR) using an NMR spectrometer in a CDCl3/deuterated hexafluoroisopropanol/deuterated trifluoroacetic anhydride mixture.\nThis measurement makes it possible to detect\nthe aromatic (terephthalic T) nuclei,\nthe oxyethylene units,\nthe mono(oxyethyleneoxy) (OEO) residues present in the oxyethylene aromatic diester (OAD) groups,\nthe polyoxyethylene terephthalate blocks (TE-POE), the units at the chain ends.\nThe calculation makes it possible to deduce therefrom:\nthe amount by weight of all the (TE) units in the terephthalic polyester composition (TPC),\nthe amount by weight of mono(oxyethyleneoxy) (OEO) residues in the terephthalic polyester composition (TPC).\nFurther information with regard to these analyses is given in the examples."}
{"text": "Subscription television systems typically include an electronic program guide (EPG). With the development of interactive digital set-top boxes (STBs), a new type of advanced multimedia EPG was developed which may be referred to sometimes as interactive program guides (IPGs). Today, the acronyms EPG and IPG have been used interchangeably. The EPG provides an onscreen display of a program grid menu that includes a program lineup and episode information for every channel for the next several days. The program grid menu may be arranged in various ways. The program grid menu itself is a user interface where the user may navigate the program grid menu using arrow keys on a remote. For example, a typical EPG would have a program grid menu based on available channel and time. The user can navigate or scroll through the list of programs within a selected time slot until the desired program is highlighted. The boundaries of other fields associated with the highlighted field may also be highlighted such as the corresponding channel and time. The highlighted field may also sometimes be referred to either as a pill or as being active. Once the desired field is highlighted as a selectable item, the user may select the highlighted field using the select button on the remote.\nTypically, a field in the EPG is shown as being highlighted by using a different color background relative to the background of any other field. Alternatively, a selectable field may be depicted as being enlarged or oversized. For the purposes of describing the present invention, the surrounding boundary or edge of the field corresponding with a selectable field is bolded to depict a field being highlighted. Although the present invention illustrates EPGs having a particular configuration or number of fields, today's EPGs include hundreds or even thousands of possibilities which the user may scroll through, highlight and select. Therefore, because of the increasing number of options within an EPG, as well as because of the difficulty some users have determining which item is selectable, there is a need to bring greater attention to the selectable item within the EPG."}
{"text": "Gas operated auto loading firearms, and particularly shotguns, typically comprise a gas cylinder parallel to the barrel of the gun. A bleed orifice extends from the barrel into the initial volume of the gas cylinder, to permit the pressurized gas resulting from the shot being discharged to pass from the barrel into the gas cylinder. The pressurized gas acts upon a mechanism which operated to replace spent shells with live shells from a magazine tube and cock the hammer of the shotgun for the next shot. For most gauges of shotgun, two or more ammunition loads are generally available. A light ammunition load is typically used for target practice, with increasingly heavy charges found in field loads and magnum loads.\nIn the past, different barrel and inertia sleeve assemblies have been used to accommodate the different ammunition loads, and effort has been directed toward the development of a firearm that would satisfactorily accommodate all of the readily available ammunition loads without modification of the mechanism or shooter adjustments."}
{"text": "The present invention relates to an optical pickup and an optical disc apparatus having a mounted optical pickup. The problems concerning wavefront aberrations such as coma aberration, astigmatism, and spherical aberration occur when beams of light are focused on an optical disc. JP-A No. 2000-40249 discloses a technique for reducing the wavefront aberrations of the light beams. In the technique of JP-A No. 2000-40249, wavefront aberrations are corrected by assigning a required phase difference to the light beams that pass through a liquid crystal element disposed in the optical path of an optical pickup.\nJP-A No. 2002-140831 discloses another technique for reducing wavefront aberrations of the light beams. In the technique of JP-A No. 2002-140831, an aberration-correcting optical system for correcting coma aberration and spherical aberration is disposed in the optical path of an optical pickup."}
{"text": "Cosmetic powders are currently available in loose and pressed form. The majority of these powders, face powder, blusher, eye shadow, etc., are used in pressed form. This form of the product is very practical since it is easy to carry along and can be used whenever needed for touching up. The desirable properties of a pressed powder are: easy powder release without dustiness, good skin adhesion, and a strong cake which does not crumble, break, or cake up. In order to achieve all of these characteristics, the present pressed powder formulations contain a blend of several ingredients, each chosen for their specific quality.\nThe major constituents of a regular shaded pressed powder are talc and color; all other ingredients being used to achieve compressibility. Binding is a special property inherent to pressed powder formulation. Most powder constituents found in formulations are not \"binding\" in themselves and therefore other agents such as metallic stearates, kaolin, and fatty materials are used as additives to insure proper adhesion and compressibility. Fatty materials, such as mineral oil or fatty alcohols, are used in liquid form to aid in binding, improve adhesion to the skin, and reduce cake dustiness.\nIn frosted products, mica and titanium dioxide coated mica (pearlescent materials) are the major powder constituents. In order to compress these materials, a large portion of about 20 to 40 percent of a composition is added. Such large amounts of non-pearlescent ingredients reduced the \"frosty\" effect of the product.\nWith current methods, manufacturing of a pressed powder is long and tedious. It involves blending of all dry powder ingredients, wetting the blend with the liquid fatty material and can lead to subsequent balling up of the powder. In order to break up the powder balls and achieve a uniform liquid distribution within the powder, it must then go through a micropulverization process. The inability to hold a shape after compression may vary with the chemical origin of the binding agent and the percentage utilized in the pressed powder formulation. A great many such formulations are possible. See for example, De Navarre, The Chemistry and Manufacture of Cosmetics, Vol. 4 Continental Press, Fla. 1974.\nIn the past many methods of aggregating inherently non-cohesive powders were used. For example, Ornfeldt, in U.S. Pat. No. 1,996,168 discloses a jelly-like substance which might be used. 1-5% of beeswax or synthetic wax and a small amount of soap are combined in melted form with mineral oil and heated. After cooling the mixture, powdery materials such as talc are added. The jelly-like substance will increase the surface cohesiveness of the powder and keep it from dusting. No mention is made of compressing this powdery product, nor is it likely that this would constitute a suitable method for aggregating a powder which is to remain substantially dry, such as face powder.\nPicker, in U.S. Pat. No. 1,707,684 employs olive oil as a waterproof binder for powder in loose form. Olive oil might be used as a cosmetic additive for its emolient properties, however, it does have a tendency to turn rancid. Moreover, it is doubtful that olive oil would provide sufficient cohesiveness to form a strong cake of powder.\nKole in U.S. Pat. No. 3,300,387 discloses a pressed powder consisting of an anti-perspirant coated with 5-15% of a water soluble wax-like material such as PEG and its methoxy derivatives mixed with a powder base and an oily binder such as mineral oil, lanolin, vegetable oils and isopropyl esters of fatty acids. The pressed powders of U.S. Pat. No. 3,800,034 to Kircher et al. contain binders such as gum, cellulose derivatives, gelatin, lignin, PVP, PVA and a complex magnesium silicate and lubricants such as metal stearates, liquid paraffin, fatty alcohols, fatty acids and oils. Harris et al., in U.S. Pat. No. 4,279,890 use metal stearates as binder and lubricant in pressed powders. Prevention of cracking in pressed powders was provided in U.S. Pat. No. 4,305,931 to Kawano et al. by incorporating a hydroxypropyl-etherified glycolipid ester.\nDespite the advances noted above, manufacturing of a pressed powder with conventional ingredients is a long involved multi-step process wherein the ingredients are handled repeatedly in order to achieve the desired properties. This explains, in part the high cost of certain cosmetics and their failure to meet performance criteria. Anti-perspirant sticks can be produced by combining active particulate ingredients such as aluminum chlorohydrate with a low melting point wax and melting the mixture so that an anti-perspirant stick can be formed from the hot wax. This is an expensive process in terms of energy consumption and results in a wet, greasy, wax-like product which is unpleasing to the touch. This process is frequently used to produce anti-perspirant sticks because a suitable compressing aid which would allow the dry manufacture of many such products is unavailable."}
{"text": "The present invention relates to a method for the manufacturing of cutting tools, such as milling cutter bodies with seats for indexable cutting inserts, and drills with seats for several cutting inserts of a hard material.\nWhen milling, the tool is given a rotating main motion at the same time as the workpiece or the tool performs a rectilinear or a rotating feeding motion. A feeding along a curve may also take place in copying milling machines and in numerically operated milling machines. For instance, the milling cutter body may consist of a face milling cutter, a corner milling cutter or a cylindrical milling cutter. These have in common that they usually carry a plurality of cutting inserts of a hard material, such as cemented carbide. These cutting inserts are normally of the indexable type, i.e. they may be detached and turned or \"indexed\" when a cutting edge has been worn out and the operator wishes to turn a new cutting edge into the operative position. To obtain a fine and smooth surface on the workpiece, it is necessary that the operative abutting edges of the cutting inserts describe substantially identical rotation orbits. For instance in a face cutting miller or a corner cutting miller the inserts should have substantially the same axial and radial positions in the milling cutter body. This implies extraordinarily high precision requirements when manufacturing the insert seats in the milling cutter body. If for instance the axial positioning is unsatisfactory, a so-called axial throw arises, which causes a deteriorated surface smoothness.\nIn several different face cutting miller applications, very high demands are set on the surface smoothness and R.sub.a values of maximally 1.5 .mu.m, R.sub.z values of between 10 and 15 .mu.m, R.sub.max of 10 .mu.m and WT values of between 5 and 8 .mu.m often occur. In order to achieve these surface criteria, it is necessary to adjust the cutting edges with a very high precision, both absolutely and relative to the other cutting edges in the milling cutter body. Thus, the axial height difference between two cutting edges should not exceed a few .mu.m. Suitably, one should be able to position cutting edges within an axial tolerance range of 5 .mu.m, and preferably even less. This has proved to be impossible because of the internal differences that occur between the positions in the milling cutter body of the insert seats, and thereby, also the cutting inserts.\nAccording to one known technique, a milling cutter body is manufactured by starting off from a piece of bar material of tool steel, such as, e.g., one with the code designation 42CrMo4H in toughened condition (hardness 270-310 HB). Shape and form are established by cutting machining, such as turning of the basic shape and milling out of the seats for the cutting inserts. Further, the bottom support surfaces of the insert seats are provided with a central, threaded hole for receiving the locking screw of the cutting insert, or a tubular shim screw for the fastening of a shim to the cutting insert. When the shaping of the milling cutter body is terminated, it is hardened. The purpose of hardening is to give the steel (in this case, normally the tool steel) a suitable hardness, either throughout its entire thickness or only to a certain depth under the surface. By heating to the austenitic range, followed by a fast cooling in water, oil or air (depending on the composition) martensite is obtained. The hardness of the martensite depends on the carbon content. Milling cutter bodies are either locally hardened in the regions around the insert seats, or they are hardened throughout, normally to a hardness degree of between 40 and 52 H.sub.R C, preferably between 43 and 47 H.sub.R C. During this hardening and the subsequent oil cooling, it is unavoidable that yieldings and distortions in the material bring about minor shifts of the positions of the insert seats. Thus, differences of 50 to 60 .mu.m may for instance occur between the axial positions of two insert seats. Thereby, the positions of the cutting inserts will of course deviate to the same extent. In many cases, this may turn out to be unacceptable, and therefore, in some applications, a possibility must exist to finely adjust the cutting edges axially, for instance by fastening the cutting inserts in movable cassettes; see for instance Swedish Patent 501 915 (Hessman et al.).\nThus, a primary object of the present invention is to form cutting positions in a milling cutter body with an increased position precision.\nA second object of the present invention is to reduce the axial throw between two insert seats in a milling cutter body to less than 20 .mu.m, preferably 10 .mu.m.\nStill another object of the present invention is to avoid complicated constructions with a plurality of different machine parts."}
{"text": "Magnetic storage devices (e.g., hard disk drives, etc.) typically use an inductive ring head to write data using magnetic fields on a recording medium. Longitudinal storage systems typically write the magnetic elements such that the polarity of the magnetic elements is parallel with the surface of the recording medium. Perpendicular recording systems typically write the data such that the polarity of the magnetic elements are perpendicular to the surface of the recording medium and therefore more dense. FIG. 1A shows an exemplary perpendicular recording system.\nPerpendicular recording system 151 comprises a write head (e.g., a monopole inductive element) having a main pole 152 and a return pole 153. When a current is passed through the coil of the main pole 152 a field 154 is created. In a typical perpendicular recording system, a soft underlayer 159 is disposed beneath recording medium 158 and creates a return path for field 154 to return to return pole 153. Depending upon the polarity of a field 154, a magnetic elements (e.g., 156 and 157) is altered on recording medium 158. For example, the polarity of field 154 causes the north pole and south pole of element 156 to be aligned in a first direction. However, changing the polarity of field 154 when writing element 157 results in an opposite polarity so that the north pole and south pole are aligned in an opposite direct on to that of element 156.\nPerpendicular storage systems facilitate greater storage density and improved bit detection and error correction characteristics over longitudinal systems. However, the soft underlayer (e.g., 159) creates problems with data erasure which were not generally encountered with longitudinal recording systems. Referring again to FIG. 1A, when a longitudinal external field (e.g., HIon 160 of FIG. 1A) is applied, soft underlayer 159 accentuates the collection of magnetic charges in the corners of structures, such as the recording head, which are in general contact with the air bearing surface. This can result in unwanted stray magnetic fields at those corners which can potentially erase stored data from recording medium 158.\nFIG. 1B shows a top view of perpendicular recording, head 151. As shown in FIG. 1B, the magnetic flux becomes more focused in the region of perpendicular recording head 151. More specifically, the flow of the magnetic flux is focused at the corners (e.g., 181, 182, 183, and 184) of perpendicular recording head 151. Accordingly, magnetic charges gather at the corners (e.g., 181, 182, 183, and 184) of the recording head and can result in corner stray magnetic fields which can write on unwanted tracks or other locations, thereby erasing data from the storage medium.\nFIG. 2 is a top view of an exemplary recording head. In FIG. 2, recording head 151 has beveled edges 190 and 191 on the corners of the surface which is exposed to the air bearing surface. It has been found that beveling the edges (e.g., beveled edges 190 and 191) of the recording head can reduce the amount of magnetic charge that can accumulate proximate to the storage medium. Instead, the build up of magnetic charges is generally in the regions of corners 190a and 191a. Furthermore, it has been found that a shallow angle of 10° or less (e.g., angle Θ of FIG. 2) is preferable to reduce the build up of magnetic charges proximate to the recording medium.\nHowever, fabricating a recording head with the cross section shown in FIG. 2 is difficult and often prone to manufacturing errors. For example, if recording head 151 is maintained at a fixed width, the shallow angle of beveled edges 190 and 191 reduces the useful area (e.g., region 197) of recording head 151 to an unacceptable dimension. Alternatively, if the useful area of region 197 is widened to an acceptable dimension, the overall width of recording head 151 may be unacceptable.\nAdditionally, when the bottom surface of recording head 151 is being defined during fabrication, precise alignment of the definition mask is required. If the definition mask is mis-aligned, too much of the bottom surface of the recording head 151 may be removed, typically indicated by line 195, thus eliminating the beveled edges. For example, an error of as little as 0.5μ can eliminate the beveled edges from the beveled recording head. As a result, the advantage of beveling edges 190 and 191 is lost and unwanted stray magnetic field effects may be exhibited by recording head 151."}
{"text": "It is known to make food products comprising a biscuit part and a filling part. For example, products having two biscuits and a filling therebetween are known as sandwich biscuits. It is also known for these food products to include in the filling-part an anhydrous filling with live lactic cultures. Anhydrous fillings comprising live lactic cultures are usually produced with yoghurt.\nYoghurt is the product of milk fermentation by symbiotic cultures of Streptococcus thermophilus and Lactobacillus delbrueckii subspecies bulgaricus. Yoghurt is known to have beneficial properties, such as facilitating lactose digestion. These beneficial properties arise from the presence of the two strains of bacteria in large quantities: in particular, the yoghurt has a total cell count of the population of both strains of over 107 cfu/g (cfu: colony-forming unit). However, the shelf-life of yoghurt is usually short; usually less than 2 months for fresh yoghurt. Therefore, yoghurt is usually kept at a low temperature to limit decay of the lactic cultures and organoleptic deterioration.\nIt is known to make yoghurt powders since they have a longer shelf-life with respect to the organoleptic properties. They are obtained via various drying methods such as freeze-drying and spray-drying. For economic reasons, spray drying is more commonly used. However, spray-drying is more detrimental to lactic culture survival than freeze-drying during production and, thus, results in a greater loss of lactic cultures during drying. Consequently, most of the commercially available yoghurt powders only contain limited amounts of live lactic cultures; lower than 107 cfu/g of yoghurt powder. These are, moreover, typically not sufficiently viable during storage.\nUS2010136182 discloses a fat-based filling containing starch.\nEP2671453 discloses a food product having a filling with a high amount of live lactic cultures.\nUS2008/305210 discloses cultures encapsulated with compound fat, as well as breakfast cereals coated with the same.\nUS2009/304864 discloses a powder of fermented milk or yoghurt with a high density of lactic ferments.\nWO99/09839 discloses a paste-like composition for use as a filling containing a significant amount of probiotic.\nWO2006064959 discloses a lactic acid bacteria-containing chewing gum tablet.\nFR2895877 provides a method for producing yoghurt powder with high amount of live lactic cultures, i.e. higher than 5×108 cfu/g of yoghurt powder. This yoghurt powder can be used in fillings for producing sandwich biscuits with a high cell count. The cell count remains high enough to be compliant with the regulatory requirements for yoghurts for up to 7 months at room temperature. However, the products are insufficiently stable to have a shelf-life of 9 months, or even 12 months, which is typical for filled biscuits. In addition, in warmer climates where the sandwich biscuits may be exposed to temperatures above 30° C. for significant periods of time, the shelf stability is unlikely to be sufficient. There is therefore a need to improve the shelf-stability of the composite product."}
{"text": "This invention relates to respirators in general and, in particular, to a new and useful respirator with an oxygen-releasing chemical cartridge."}
{"text": "The manufacture of ceramic microelectronic packages by the co-firing process involves deposition of trace refractory metalization on the unfired, or \"green\" tape cast of the ceramic IC carrier. The trace metal represents the package bond pads which the IC die leads attach to. The firing of the green tape shrinks the ceramic to a final, post-fired size. The shrinkage is predictable (on the order of 18%) to within a +/- \"acceptable variation\" tolerance. The standard variation tolerance is +/-1.0% of the nominal post-fired ceramic body dimension.\nThe trace metal being in the ceramic also shrinks predictably when the ceramic is fired, with the same +/-1.0% variation tolerance on size and geometrical spacing. If a first array of trace line metal is deposited on the green tape ceramic and is intended for post-fired registration with a second array of trace metal having a nominal dimension of 1.0 inches, then the final dimension of the first array may be 1.000+/- 0.010 inches.\nPrior art methods of mounting the IC dies to the ceramic package bond pads include both manual and automated bonding processes. In the manual processes the electrical interconnection of the IC to the package require operator alignment of the IC interconnect wires (or ribbon) to the bond pads. After alignment a bonder machine secures the device leads to the bond pads by welding or soldering. Automated processes, such as tape automated bonding (TAB), require accurate positioning of the package bond pads so as to automatically receive the IC TAB leads in register. The +/-1.0% shrinkage tolerance creates problems with the TAB leads to bond pad registration.\nWhile the bond pads may be made large enough to accommodate all manufacturing tolerances, including the shrinkage tolerance, this prevents increasing package density. The IC industry drive towards miniaturization requires high density packaging which prohibits this brute force method of lead registration."}
{"text": "Chat applications and/or systems may facilitate synchronous textual communications between two or more users. Graphical user interfaces presenting chat sessions may present the textual communications sent by each user in real-time and/or near real-time. Typical graphical user interfaces present chat sessions having one or more different participants as separate windows and/or graphical user interfaces.\nGroup chat sessions may facilitate synchronous textual communications between a group of users. New groups with one or more different users may be created such that users often participate in multiple chat sessions with different individual other users and/or groups of users."}
{"text": "1. Field of the Invention\nThe present invention relates generally to object sensors and related methods, and in particular relates to electronic object sensors and methods useful in detecting vehicle speed and shape for classification and input to Intelligent Vehicle Highway Systems (IVHS).\n2. Background Art\nA vehicle sensor detecting the presence of a vehicle in a traffic lane and indicating the vehicle speed as it passed the sensor is described in U.S. Pat. No. 5,321,490 referenced above. A time-of-flight laser rangefinder is used to measure the normal distance to a road surface from a fixed point above the road surface and then measures the distance to a vehicle which either passes or stops under the sensor. Two laser beams pulsing at a high rate are projected across the road surface at a fixed angle between them. Because of the high repetition rate of the pulsed beam, the system is able to determine vehicle speed with an accuracy within one mph and using this calculated speed, develop a longitudinal profile of the vehicle using consecutive range measurements as the vehicle moves under the sensor. Such active near-field object sensors are relatively low in cost, are accurate and have utility in a wide variety of applications. A laser diode capable of emitting pulses of coherent infrared radiation is used together with collimating optics and a beam splitter to provide two diverging output beams directed toward the road surface under observation.\nThe sensor receives a portion of the energy reflected from either the area, or an object located within the area, such as a vehicle. The returned pulse energy is then provided as an input to a receiver for determining a time of flight change for pulses emitted and received, which may be caused by the presence of an object within the area. The sensor is also provided with various features useful in providing outputs which indicate either the speed, census, size or shape of one or more objects in the area. For example, the sensor is provided with means for receiving an input from the time of flight determining means and for providing an output indicating whether the object meets one of a plurality of classification criteria (e.g., is the object an automobile, truck or motorcycle). A receiver includes two detectors and alternately selects between outputs of the two detectors for providing the time-of-flight measurements. The time interval between interceptions of the two diverging outputs by a vehicle provides the speed of the vehicle passing through the area.\nU.S. Pat. No. 5,278,423 referenced above discloses the generation of three-dimensional images of objects by rotating or scanning a laser beam rangefinder, operating at a high pulse rate, in a plane where there is relative motion between the rangefinder and the object to be sensed or imaged in a direction perpendicular to the laser beam plane of rotation. The laser rangefinder rotating beam covers the object being sensed permits a three-dimensional image of the object to be determined. By way of example, the '423 patent discloses a sensor traveling between rows of trees with the laser rangefinder scanning on either side of a moving vehicle carrying the sensor. Beam scanning is within a plane perpendicular to the motion of the vehicle. When the sensor detects the presence of foliage, it provides a signal activating a spraying system for the efficient spraying of the tree. This operation ensures that spraying takes place only when there is foliage present to intercept the sprayed materials. Economic and environmental benefit is thus realized.\nThe agricultural sprayer employs a pulsed time-of-flight range measuring system having separate apertures for a laser transmitter and receiver. The laser beam and receiver field-of-view are continuously scanned by a rotating mirror in a vertical plane which is perpendicular to the forward motion axis of the sprayer vehicle. The position of the mirror, and correspondingly the laser beam, is determined by a shaft encoder attached to the mirror drive motor shaft. With this embodiment, a single sensor makes range measurements on both sides of the vehicle as the vehicle moves the sensor between rows of trees. Since the sensor only needs to detect the presence of trees, range measurements are only made within elevation angles of plus and minus 45 degrees on each side of the sensor. Data is collected within 180 degrees out of the 360 degrees of a revolution or circular scan. As the vehicle moves along, the scan trace advances on consecutive revolutions of the mirror. Employing a distance traveled input from the vehicle, the sensor creates a panorama of images. An algorithm then determines whether trees are present from the measured range data as a function angle. Spray units are grouped in zones and the sensor provides foliage images for the zones and thus an indication of the amount of spray necessary for a particular tree zone.\nThere is a continuing demand for accurate, low-cost sensors useful in a wide variety of applications, including equipment used in the home, as well as for security, military and transportation applications. Traffic signal controllers utilizing overhead sensors are known as described by way of example in U.S. Pat. No. 3,167,739 to Girard et al; U.S. Pat. No. 3,436,540 to Lamorlett; U.S. Pat. No. 3,516,056 to Matthews; U.S. Pat. No. 3,532,886 to Kruger et al; U.S. Pat. No. 3,680,047 to Perlman; and U.S. Pat. No. 4,317,117 to Chasek. Likewise referenced, near-field sensors have also been utilized as intruder alarms and as automatic door operators. Examples of such arrangements are disclosed in U.S. Pat. No. 3,605,082 to Matthews; U.S. Pat. No. 3,644;917 to Perlman; U.S. Pat. No. 3,719,938 to Perlman; U.S. Pat. No. 3,852,592 to Scoville et al; U.S. Pat. No. 3,972,021 to Leitz et al; and U.S. Pat. No. 4,433,328 to Saphir et al. U.S. Pat. No. 4,768,713 discloses the use of an ultrasonic sonar-type sensor to detect the presence of tree foliage, as do U.S. Pat. Nos. 4,823,268 and 5,172,861. Optical dimensioning techniques have been incorporated in industrial uses as disclosed in U.S. Pat. Nos. 4,179,216 and 4,490,038.\nVehicle detection and classification provided by the Intelligent Vehicle Highway System and Method of the above-referenced related inventions, and as herein described, have proven to be very successful using rule-based algorithms to do shape-based vehicle classification. With the current class structure, the sensors have achieved 98.5% accuracy on a random 10,000 vehicle test. Using the National Academy of Sciences (NAS) vehicle database, which was collected under the ITS-6 IDEA program, the sensors achieved 96.5% accuracy on a 50,000 vehicle database including a range of weather conditions and traffic conditions.\nHowever, in spite of this success, most applications in the U.S. still require vehicle classifications based on the number of axles a vehicle has as opposed to just the shape-based classification. This implies that tolling applications in the U.S. must deploy an axle counter, such as a treadle, along with a host of other sensors to perform the other tasks that the axle counter cannot do, such as vehicle separation. Along with the disadvantages of embedding a sensor in the road, the higher system cost of using several sensors, and added complexity of integrating different sensors into a single system, transportation professionals at all levels have expressed an increased interest in an overhead sensor that will count axles."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and arrangement for preparing a neural model of a paper winding device, such as a drum cutter.\n2. Description of the Related Art\nGiven the production of paper, it has proven successful to store the paper in the form of drums, which are also referred to as tambour or coil, since very large paper quantities thereby can be stored crease-free in the smallest space for bearing, transporting and processing.\nA paper winding device contained in the drum cutter has the task to unwind a broad coil of paper and to develop paper webs corresponding to the requests of the customers, whereby conventional paper widths can certainly amount to 9.5 m. For this purpose, these are cut in longitudinal direction and the narrower paper webs arising thereby are wound up again. Doubling roll coilers are frequently utilized due to their technical preferences for the winding-up (one is schematically shown in FIG. 1). This winding process critically influences the quality of the paper on the wound rolls, this quality being strongly dependent on the tensions in the paper occurring in the rolls. The tangential and radial tensions arising in the paper coil, among other things, are influenced by the following influencing and control quantities of the paper winding device and the paper:\nby their moments and speed given the center drive;\nby its moment and speed given the doubling roll drive;\nby the winding radius, the line force in the nip, adjustments of the friction dampers and the width of the coils given the winding stations;\nby geometrical influencing quantities, such as the doubling roll radius, the width of the tambour, and the radius of the sleeve onto which the paper coils are wound;\nby the properties of the paper, such as the modulus of elasticity, the G.M.S per density, the roughness, the smoothness, the moisture, the porosity, as well as by the elongation at rupture;\nby other influencing quantities, such as the web strength, the roughness of the doubling roll, the friction value of the doubling roll for each paper type, as well as by the modulus of elasticity of the sleeve onto which it is wound.\nSince these influencing and controlling quantities are also dependent on environmental conditions, such as the humidity and the temperature, such winding processes are very difficult to describe analytically. For these reasons, it is particularly difficult to meet the quality requirements with respect to such coils and to constantly keep them. Expert knowledge is still utilized thereby with respect to the prior art, in particular Predominantly, it is to be assured that reproducible coils having an optimal coil hardness structure arise during the winding process. Furthermore, a lateral progressing, which is also referred to as telescoping of the roll, is to be prevented. The winding process is to be particularly rid off tears and broken picks, as well as plastic deformations.\nThe controlling quantities for such paper winding devices and particularly for drum cutters have been hitherto fixed, in the form of set value tragjectories in an initial operation phase by initial operation engineers, manually and on the basis of their expert knowledge. This requires great technological knowhow, skills and time. Generally, a separate set of set value trajectories of the paper winding devices is fixed for each paper type of the product spectrum and is stored in the, for example, process computer, which controls the drum cutter. In a later operation, the set value trajectories are read out of the memory and are utilized for controlling the winding device. A specific quality level is achieved as a result of this course of action; however, the quality that can be obtained thereby is limited by the following factors:\nit cannot be assured that optimal control trajectories are always utilized, since there is no mathematical model for determining optimal control trajectories. On a case by case basis, sub-optimal control trajectories therefore can be applied.\ndrifts of measuring devices or as a result of a tear remain unconsidered.\nvariations of the paper parameters within one type, such as the G.S.M or the elasticity module, remain unconsidered.\nWhen the winding quality worsens below an acceptable value as a result of the cited effects during the operation of a drum cutter or of another paper winding device, the driver of the device must manually intervene for correction purposes. General theoretical basics for the processes occurring during the paper winding are described in H.-J. Schaffrath, F. Hibinger and L. Goettsching xe2x80x9cNumerische Berechnung von Spannungsverlaeufen und Wickelhaerte in einer Papierrollexe2x80x9d, pages 350-361, No. 6, 1994. Methods for the automated preparation of control parameters for paper winding devices and particularly for drum cutters are not yet known.\nThe present invention is based on the object of providing a method and a device for the neural modeling of paper winding devices and particularly for the optimization of the control parameters of these devices.\nThis and other objects and advantages of the invention are provided by a method for the neural modeling of a paper winding device, wherein influencing and control quantities of a paper winding device are provided in a first step for the training of a neural network, in that these are measured and are stored at least depending on the wound number of layers and the appertaining coil radius of the paper coil, wherein the web strength in the paper web or a quantity correlated with it, depending on the influencing and control quantities of the paper winding device, is determined as target quantity for the neural network from the relationship between the measuring quantities coil radius and appertaining number of layers of the paper coil, and wherein the neural network is trained as model of the paper winding device via a conventional learning method at least with influencing and control quantities of the paper winding device as input quantities and the web strength dependent thereon or, respectively, the quantity correlated with the web strength as output quantities.\nIn the method, the influencing and control quantities are measured at a real paper winding device and/or at the paper, and/or the environment of the paper winding device or of the paper. The influencing and control quantities are measured dependent on time.\nIn one embodiment, the method for generating optimal control quantities for a paper winding device depending on a desired winding quality of the paper coil, wherein influencing quantities and control quantities of the paper winding device to be optimized are supplied in a first optimization step to the neural model prepared according to one of the claims 1 through 3 and wherein a first web strength or, respectively, a first quantity correlated with the first web strength is determined therefrom with the aid of the model, wherein a first winding quality is determined in a second step from the known relationship between the web strength or, respectively, the quantity correlated with the web strength and winding quality from the first web strength or, respectively, the correlated first quantity, and wherein the first winding quality is compared to the desired winding quality and wherein at least second control quantities of the paper winding device to be optimized are supplied in a further optimization step to the model, whereby this course of action is kept until the winding quality determined from the two steps above sufficiently accurately corresponds to the desired winding quality.\nThe change of the winding quality is determined depending on the change of the control quantities of two adjacent optimization steps and wherein it is concluded from this change in which way the control quantities are to be changed for the next optimization step in order to arrive at the desired winding quality. The control trajectories may be generated as a time-dependent sequence of control variables.\nAn operating method for a paper winding device is provided, wherein control quantities, which have been generated in an optimized manner according to the foregoing, are supplied to the paper winding device.\nThe objects of the invention are also achieved by an arrangement for the neural modeling of a paper winding device, which includes first means for measuring influencing and control quantities of the paper winding device depending on the wound number of layers and the appertaining coil radius of the paper coil, which includes second means for determining the web strength in the paper web as a target quantity or, respectively, a quantity correlated with the web strength for the neural network at least from the relationship from the quantities coil radius and appertaining number of layers of the paper coil measured by the first means, and depending on the influencing and control quantities of the paper winding device, whereby the quantities measured by the first means being supplied to said influencing and control quantities, which has a neural network as a model of the paper winding device, whereby the quantities, which are determined or, respectively, measured by the first and second means being supplied to said neural network, which learns via a conventional learning method the dependency of the quantities from one another, which are supplied to it.\nPreferably, the arrangement at least includes means for storing one of the measured or determined quantities. The paper winding device may be fashioned as a tambour cutter.\nA method provides at least driving moments of the winding device, as well as the elasticity module of the paper are measured.\nA particular advantage of the inventive method is that the neural model is prepared with the aid of an intermediate quantity, the web strength or a measured quantity correlated therewith, whereby it can be determined from the relation between the coil radius and the number of layers. Since the measuring quantities are determined at the real system for this relationship and are available for different points in time, a neural network can be directly trained as a model of a nip by means of control parameters. The advantage is that merely the relationship between control quantities and web strength or a quantity correlated with it must be modeled by the neural network and that relationships that are known from the prior art can be used for determining the winding quality depending on the web strength. When the average layer thickness is used as a quantity correlated with the web strength, it can be determined by measuring the radius increase given counting of the number of layers at the same time. The cited prior art in the mentioned printed publication is referenced for detail information regarding this topic.\nAdvantageously, the proposed method can be applied in a two-step fashion. For example, while a neural network is coupled to a paper winding device and the influencing and control quantities are supplied to it, it can learn the relationship between the number of layers and the control parameters, as well as of the web strength or, respectively, the quantity correlated with it. A second neural network can be utilized as a current copy of the neural network connected to the machine for preparing optimized control parameters. Advantageously, an iterative method is carried out for determining the optimized control parameters, whereby the relation between the control and influencing quantities and the web strength or, respectively, the quantity correlated with it is modeled particularly by the neural network, which represents a model of the paper winding device and thereby particularly of the nip, and whereby the winding quality is calculated with the modeled web strength from the known relationship from the prior art. When a quantity correlated with the web strength is modeled, it is possible that it is directly measurable, as this is valid for the average layer thickness. This calculated winding quality can be subsequently compared to the desired winding quality to be obtained during the winding process, and the control parameters can be adapted to the operation, in a quasi off-line manner, solely by utilizing the model in that they are continuously modified until the desired winding quality appears with sufficient accuracy.\nAn optimal changing direction of the control parameters can be concluded from the difference of the winding quality between two adjacent optimization steps and changes of the control parameters, which are associated therewith and which relate to these optimization steps, so that control parameters, which are changed in the correct direction, can be prescribed in a further optimization step.\nAs a result of the proposed method, set value trajectories are advantageously provided for the control parameters, since the measured quantities change in the course of the winding process and adapted control quantities are required.\nA paper winding device or a drum cutter can be particularly advantageously operated by means of the optimized control parameters or, respectively, set value trajectories as proposed, since the above-described disadvantages are thus avoided and a higher quality level can be obtained during the paper winding process.\nThe proposed method is particularly advantageously performed with the aid of an arrangement, wherein the neural network is directly coupled to the paper winding device and wherein the relevant measured quantities, such as the influencing quantities of the paper, can be directly supplied to it, since the neural network can thus be permanently retrained in the online mode during the winding process. A copy of the neural network can be simply made for an optimization of control parameters occurring on a case by case basis for further winding processes. Therefore, the most current and best trained version of the network is always available.\nParticularly advantageously, a neural network can be arranged at the paper winding device of a drum cutter, since drum cutters are frequently utilized in order to carry out customer customizations of paper webs.\nIn the proposed method or, respectively, in the proposed arrangement, at least the driving moments of the winding device and the modulus of elasticity of the paper are particularly advantageously measured, since these influencing quantities are extremely important for the obtainable quality during the winding process."}
{"text": "Toy blowguns typically fire projectiles such as darts. They are relatively simple devices, but have always been popular as toys or for use in recreational activities. A user simply inserts a projectile inside a pipe portion of a blowgun and then blows a powerful, short blast of air into a mouthpiece at one end of the pipe portion. The force of the blast of air provides propulsive power to the projectile, launching it from the pipe into the air towards an intended target.\nThe propulsive power depends on the respiratory power produced by the user. As such, the use of blowguns provides exercise benefits, especially in developing the lungs.\nChildren, however, typically do not have sufficient respiratory power to use blowguns, and even toy blowguns, effectively. This is especially so for younger children. In particular, it is difficult for children to fire projectiles from prior blowguns at speeds or at distances that make the use of these blowguns particularly fun or interesting.\nIt is an object of the present invention to overcome or ameliorate at least one of the disadvantages of the prior art, or to provide a useful alternative."}
{"text": "In order to print data, such as text, the data to be printed is encoded. The encoding schemes represent the characters to be printed. Because there are numerous written languages having distinct characters, the encoding schemes represent, for example, the characters in different languages. Although there are various schemes utilized to encode data, Unicode is a standardized code for representing text. Unicode includes a code word for each character in a particular language. For example, one code word in Unicode would represent a Chinese pictograph, another code word in Unicode represents an Arabic character, while yet another code word in Unicode represents a character in the Latin alphabet. The code words also have a particular length. Unicode uses a two byte code word for most characters. In addition, Unicode groups characters in a particular numeric range. For example, all Arabic characters are represented by Unicode code words within a particular numeric range.\nIn addition to using a particular encoding scheme, a font is also used to print the data stream. For example, an Open Type font might be used in conjunction with Unicode. The font specifies the glyph, which is the character shape actually printed by a conventional printer. Thus, the font would indicate whether the “A” represented by the Unicode code word is a Times New Roman glyph or an Arial glyph. The font could also include Chinese pictographs, Arabic characters or other characters. The font is used to convert between the code word and the glyph. For example, the font might include an encoding table and a glyph table. The encoding table would be used to convert the code word for a particular encoding scheme to a corresponding glyph index. The glyph table would be used to convert the glyph index to a particular glyph. The glyph would then be printed in a particular location on the page corresponding to the code word.\nUsing Unicode code words in a data stream, a conventional computer system could convert Unicode to other data, for example a bit map of the glyph to be printed and the location of the glyph. This data is provided to a conventional printer. The conventional printer would print the glyph in the location corresponding to the Unicode code word. In such rendering a particular code word or set of code words always corresponds to the same glyph or set of glyphs. Thus, one-to-one rendering of a code word to a glyph is performed.\nAlthough encoding schemes such as Unicode provide a code word for each character, one of ordinary skill in the art will readily recognize that in certain, complex languages, the glyph used to represent a given character depends upon its context. For example, in Arabic, Indic, and Thai language groups, the glyph used depends upon the surrounding characters. In addition, certain languages are written in a different order. English is typically written from left to right. However, other scripts are written from right to left. In addition, some scripts are written both right to left and left to right. For example, in Hebrew and Arabic, text is written right to left, but numbers are written left to right. As a result, the processing of code words is highly dependent upon the language. For complex languages, for example those read in a different order or having characters that are context dependent, identity and/or order of glyphs depend upon the context. Thus, printing of complex text data, such as the complex text described above, cannot be performed using one-to-one rendering.\nFIG. 1 depicts a conventional system 10 used to print data including complex text data. The conventional system 10 might print text using an encoding scheme such as Unicode. For clarity, the conventional system 10 is described in the context of Unicode. The conventional system 10 includes a computer system 20 and a printer 30. The computer system 10 might include one or more clients and a print server or a desktop computer. The conventional computer system 20 includes a conventional formatter 22 and, for certain complex text data, a conventional layout engine 24, and a conventional font 26. The conventional formatter 22, conventional layout engine 24, and font 26 might reside on a print server (not shown) or a desktop. The conventional formatter 22 formats a data stream to be printed. The data stream is encoded in Unicode. Consequently, the conventional formatter 22 uses the conventional layout engine 24 for Unicode complex text. The conventional layout engine 24 might be a Unicode layout engine such as the Uniscribe layout engine for WINDOWS. The conventional layout engine 24 can typically run on platforms such as Windows or UNIX. The font 26 may be an Open Type font that is used to render complex text.\nThe conventional layout engine 24 analyzes complex text data, such as Unicode complex text, and determines the appropriate glyph for a particular Unicode code word. The conventional layout engine 24 thus includes data related to languages and contexts. The conventional layout engine 24 utilizes this data to analyze complex text, determine contexts and thus determine the appropriate glyph. The conventional layout engine 24 thus also uses the font 26. For example, if the font 26 is an Open Type font, the conventional layout engine 24 would determine the appropriate glyph and, using the font 26, determine the appropriate glyph indices. Alternatively, the conventional layout engine 24 might determine a bit pattern for the glyph using the conventional rasterizer 28. This would occur when the printer 30 does not support the font 26. For example, the conventional layout engine 24 might use the font 26 to determine the appropriate glyph indices and, using the glyph indices, the appropriate glyph. The conventional rasterizer 28 might then output a bit pattern of the appropriate glyph.\nFIG. 2 depicts a conventional method 50 for printing text encoded using Unicode and including Unicode complex text. The print data is obtained, for example from a client, via step 52. The conventional layout engine 24 in the computer system 20 is used to convert any Unicode complex text in the print stream into glyphs, via step 54. In order to perform step 54, the conventional layout engine 24 utilizes the font, described above. In converting the Unicode complex text into glyphs in step 54, the conventional layout engine 24 and conventional rasterizer 28 typically output the bit patterns of the glyphs.\nThe remaining portion of the Unicode text, which is not Unicode complex text, is rendered using one-to-one rendering, via step 56. The result of one-to-one rendering is bit patterns. In addition, in contrast to the Unicode complex text, one-to-one rendering might take place in the conventional printer 30 rather than in the conventional computer system 20. The conventional data stream, which includes bit patterns for Unicode complex text and Unicode code words for the remaining Unicode text, is passed to the conventional printer 30 for printing, via step 58. The conventional data stream also includes the locations of the glyphs. The conventional printer 30 then prints the data stream using the conventional rasterizer 34, via step 60. Step 60 might include performing one-to-one rendering for Unicode data that is not complex. Step 60 also includes placing the glyphs in the appropriate positions on the page.\nAlthough the conventional system 10 and conventional method 50 function, one of ordinary skill in the art will readily recognize that there are drawbacks. The data for the Unicode complex text that is passed to the printer 30 is typically in the form of bit patterns. In any case, the code words for the Unicode complex text are not passed to the conventional printer 30. Once the bit patterns are generated, neither the computer system 20 nor the conventional printer 30 has any knowledge of the text represented by the data corresponding to the Unicode complex text. The data stream can no longer be processed for text operations, such as sorting, searching, indexing, merging or repurposing. In addition, bit patterns for a glyph are relatively large in comparison to the corresponding Unicode code words. Consequently, the transmission of data between the conventional computer system 20 and the conventional printer 30 may be inefficient and time consuming.\nAccordingly, what is needed is an improved system and method for printing Unicode complex text data. The present invention addresses such a need."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of protein processing and protein signaling pathways and specifically to two novel proteins having distinct activities, which are derived from a common hedgehog protein precursor.\n2. Description of the Related Art\nEmbryologists have long performed experimental manipulations that reveal the striking abilities of certain structures in vertebrate embryos to impose pattern upon surrounding tissues. Speculation on the mechanisms underlying these patterning effects usually centers on the secretion of signaling molecule that elicits an appropriate response from the tissues begin patterned. More recent work aimed at the identification of such signaling molecules implicates secreted proteins encoded by individual members of a small number of gene families. One such family of proteins which may have an influential effect upon patterning activities are those proteins encoded by the hedgehog gene family.\nThe hedgehog (hh) gene was initially identified based on its requirement for normal segmental patterning in Drosophila (Nusslein-Volhard, C. & Wieschaus, E, Nature 287:795-801, 1980). Its functions include local signaling to coordinate the identities of adjacent cells within early embryonic segments (Hooper, J. E., & Scott, M. P. Early Embryonic Development of Animals, pp. 1-48, 1992) and a later function in cuticle patterning that extends across many cell diameters (Heernskerk, J. & DiNardo, S., Cell, 76:449-460, 1994). The hh gene also functions in the patterning of imaginal precursors of adult structures, including the appendages and the eye (Mohler, J. Genetics, 120:1061-1072, 1988; Ma, et al, Cell, 75:927-938, 1993; Heberlein, et al, Cell, 75:913-926, 1993; Tabata, T. & Kornberg, T. D., Cell, 76:89-102, 1992; Basler, K. & Struhl, G., Nature, 368:208-214, 1994). Genetic and molecular evidence indicates that hedgehog proteins are secreted and function in extracellular signaling (Mohler, J., supra; Lee, et al, Cell, 71:33-50, 1992; Taylor, et al, Mech. Dev., 42:89-96, 1993).\nIn vertebrates activities encoded by hh homologues have been implicated in anterior/posterior patterning of the limb (Riddle, et al, Cell, 75:1401-1416, 1993; Chang, et al, Development, 120:3339, 1994), and in dorsal/ventral patterning of the neural tube (Echelard, et al, Cell, 75:1417-1430, 1993; Krauss, et al. Cell, 75:1431-1444, 1993; Roelink, et al. Cell, 76:761-775, 1994).\nThe vertebrate ventral midbrain contains neurons whose degeneration or abnormal function are linked to a number of diseases, including Parkinson's disease and schizophrenia. It is known that motor neurons develop in close proximity to the floor plate in the ventral midbrain. Midbrain projections to the striatum are involved in the control of voluntary movement (Bjorklund and Lindvall, In: Handbook of Chemical Neuroanatomy, eds., Borklund, et al., Amsterdam: Elsevier, pp 55-122, 1984) and loss of these neurons results in the motor disorders of Parkinson's disease (Hirsch, et al, Nature, 334:345, 1988). Midbrain dopaminergic neurons that innervate limbic structures and the cortex influence emotional and cognitive behavior, respectively, and abnormal function of these neurons has been associated with schizophrenia and drug addiction (Seeman, et al, Nature, 365:441, 1993).\nWhile the molecular nature of the factors that specify neuronal cell fate have not been established, members of the transforming growth factor-β (TGF-β) (Lyons, et al, Trends in Genetics, 7:408, 1991) or the hedgehog protein family (Smith, J. C., Cell, 76:193, 1994) may possess the characteristics expected from such factors as they participate in specification of cell fate, mediate inductive interactions between tissues, and in many cases act at a distance of only a few cell diameters.\nThe present invention establishes that hh activities encoded by these genes play a crucial role in early patterning of the developing eye and in patterning of the brain. For the first time, the invention shows that internal cleavage of hedgehog protein product is critical for full function, and that the two novel products of this auto-proteolytic cleavage display distinguishable activities, thus demonstrating that hh signaling activity is a composite effect of two separate signaling proteins that derive from a common hh protein precursor. In so doing, the invention provides the means for specific patterning and proliferation of desired neuronal cell types for addressing disorders which arise from neuronal degeneration or abnormal function."}
{"text": "This invention relates to a system to emulate underwater mines. More particularly, the system of this invention provides high frequency acoustic responses to more thoroughly emulate acoustic sensor capabilities of some underwater mines.\nShips and other marine craft that ply the waterways radiate acoustic signals that are characteristic of their design, speed, on board machinery, etc. A specific composition of the acoustic signals is relied on in some mines to selectively initiate detonation when such signals are detected within a predetermined proximity. It is, therefore, of interest to be able to determine if such compositions of acoustic signals are being emitted in sufficient strength to trigger some mines.\nA contemporary system for detecting, analyzing, and providing information regarding the compositions of such acoustic signals is limited in bandwidth. Consequently, it is not capable of emulating high frequency acoustic mines since its hardware is limited to 10 kHz and is not capable of receiving acoustic signals between 20 to 40 kHz. In addition, unlike some threat mines, it is also not capable of comparing acoustic signals from an omni-directional sensor and a directional sensor in this frequency range. These limitations have greatly reduced its ability to emulate mines that are responsive to higher frequencies.\nThus, in accordance with this inventive concept, a need has been recognized in the state of the art for an effective means to extend the acoustic frequency response of an underwater system that senses remotely originating acoustic signals, collects acoustic data, and emulates mines.\nAn object of the invention is to provide an acoustic float system to extend the frequency response of an exercise mine system.\nAnother object is to provide an acoustic system extending acoustic frequency response to permit emulation of mines having extended frequency responses.\nAnother object is to provide an acoustic system for an exercise mine system creating a mine-fire capability that is a function of omni-directional and directional sensor outputs.\nAnother object is to provide an acoustic system having an inclinometer for measuring tilt attributed to ambient currents to aid in mine fire analyses.\nAnother object is to provide an acoustic system to extend the range of frequency response and account for inclination attributed to currents.\nAnother object is to provide an acoustic system being readily capable of being modified to emulate omni-directional and directional acoustic responses of different mines.\nThese and other objects of the invention will become more readily apparent from the ensuing specification when taken in conjunction with the appended claims.\nAccordingly, the invention extends the high-frequency response of a mine emulation system. A watertight housing has a first acoustic sensor extending through a cover portion to provide an omni-directional response to higher frequency acoustic signals impinging through ambient water. A second acoustic sensor extends through the cover portion of the housing adjacent to the first acoustic sensor. The second acoustic sensor provides a directional response to the higher frequency acoustic signals. Processing circuits in the housing receive signals from the first and second acoustic sensors and generate signals representative of the higher frequency omni-directional response and the directional response. A cable extends to a mine emulation system to feed the representative signals thereto for additional processing of the higher frequency responses."}
{"text": "In the foodservice industry, it is common to serve prepared food from steam or “Sterno” flame tables. These are devices which hold hot for serving standard sizes of “pans” of food that are typically found in canteens, banquet halls, delis, corner shops and cafeteria serving lines. Typically, the food is prepared hot in the back room and is then placed into the pans for holding hot for serving. Serious food safety problems and in particular drying out of the food, are frequently reported with those heating tables. Some industrial companies also provide similarly prepared food products as quality and convenience alternatives for foodservice operators. These products are usually sold frozen or refrigerated, and heating is required before serving to a hot temperature, i.e. at least 71° C., for safety reasons, and for consumption at temperatures where consumers can enjoy the food the most, i.e. about 55° C. These products are most commonly manufactured and sold in partial “pan” sizes (e.g., ½ and ⅓ “pan” sizes) and most commonly packaged in a lidded container such as in gauge aluminum foil.\nMost often, these products are heated to serving temperature in a conventional oven similar to a residential oven. However, this method of heating has proved to be slow, as radiation heat transfer from oven wall to the package surface is relatively poor in the case of the aluminum packaging. The aluminum package has a tendency to reflect heat, resulting in the primary mode of heat transfer to be natural convection in this case. Increase of the heat transfer may be obtained by circulating hot air around the packaging as in forced air convection ovens, which are used by foodservice operators. Another enhancement is to add steam to the air to get more rapid heating but, this method may adversely affect the quality of the food product. Other methods consist of microwave heating the food product by removing the food product from the original aluminum packaging and by placing it into a special plastic tray adapted to microwave heating. Apart from the fact this requires operator's manipulations, microwave heating provides uneven heating with hot and cold spots that require rest time for the temperature to even out. The use of a combination of microwave and convection is known which provides improved results, but is relatively expensive and requires high amperage electrical connections and large amounts of electricity and also requires special packaging adapted for microwave use. In all cases of these various types of ovens, the equipment is relatively expensive, complicated and large. Once located the equipment must be used in its “place”, most typically what is referred to as “back of the house”.\nRegardless whether a conventional or convection oven is used, often the food at the edges and corners are overheated and sometimes burnt, lowering the overall quality of the prepared food. This occurs because the food product at the edges and corners is getting heated from 2 or 3 sides and at the same time is close to the surfaces. Thus, it would be desirable to have a heating method whereby the amount of heat can be controlled to a lower amount in the areas of the tray which are prone to overheating and burning. Also, if different foods are placed in a tray (e.g., meat loaf and mashed potatoes), it would be advantageous to have different amounts of heat to be applied to different parts of the tray for various food components. In conventional or convection ovens, there is no way, for the most part, to lessen the amount of heat to these edges and corners or to control different amounts of heat to different regions of trays, except to change to more complicated and expensive packaging.\nIn the foodservice area, the food may usually come in large-size frozen portions of several kilograms. As a matter of example, a standard “half-pan” lasagna packaged in an aluminum tray of 2.7 kg requires about 1950 kJ of energy (including moisture losses) and may take more than 125 minutes in a conventional oven to be heated from frozen to an acceptable hot temperature for serving. Foodservice operators are less and less inclined to accept so long heating times. One possibility is to divide the food in several smaller portions which are put in smaller containers, trays or dishes. However, this is labor intensive and it requires more attention from skilled operators. Therefore, there is a need for conveniently, rapidly and cleanly heating food originally packed in medium to large size packages without the requirement to remove the food from their packages so that labor is kept to a minimum and the food can be easily transferred to the point of service in hot conditions.\nUsually, once the food is heated in the “back room” to serving temperature, it is brought to the “front of the house” and placed in a steam table, or a chafing dish with a “Sterno” flame, to be held hot. Steam tables, however, are quite large, not inexpensive, typically take about 45 minutes to one hour to preheat, and are very labor intensive to clean. Steam tables work by heating water to create a hot, humid air environment under the food tray. Consequently, the packaging of the food product is limited to things such as aluminum or plastic unless a special liner is used. Therefore, it would be desirable to provide equipment that is quicker, simpler and reasonably low in cost, small and portable in some situations, and easy to use.\nFood warmers are known which are designed to receive for heating an open tray-shaped member for different food containers such as dishes or pans. However, the heating capacity of those warming devices is limited as the energetic loss to the environment is high due to the open configuration. Increasing the electric power of those devices to compensate for heat loss may lead to the violation of domestic electrical safety regulations which require precise electrical power limits not to be exceeded. Thus, the food warmers are generally used just for holding food warm or heating small size portions from ambient to warm but they are not capable of handling the heating of large portions of food, especially, when the food is originally in frozen or chilled state in the plate. For instance, U.S. Pat. No. 3,043,943 to J. R. Moot relates to a portable food warmer, and more particularly to a tray adapted to heat food and dishes containing food to serving temperature. The food warmer comprises a dish supporting member, an inner heat diffusing member and heating elements between the heat diffusing member and a heat insulator. Sloping top flanges are provided to delimit the contour of the open upper cavity so that dishes of different sizes and capacity can be heated in contact to the dish supporting member.\nNumerous patents relate to foodservice systems wherein individual portions of precooked food are supported on trays installed within moveable carts. The carts are usually placed in a chilled environment to store the food in a chilled state. Selected foods on the trays are rethermalized in the cart, i.e., heated to a hot serving temperature. The rethermalization method is usually directed to conductively applying heat to the selected food portions through a heating plate which is heated by a thermostatically controlled electrical heater. Examples of those systems are described in U.S. Pat. Nos. 4,068,115; 4,235,282; 4,584,466 and 5,285,051. The trays are intended to be served to the consumer as such and therefore are not adapted to heat large size portions of food but only individually portioned food in dish.\nU.S. Pat. No. 3,608,627 relates to a combination refrigeration and cooking device. A plurality of casseroles are provided into a freezing or chilling chamber. Each casserole has an electric heating element associated therewith which is connected to an electric circuit controlled by suitable switching and/or timing means. Each casserole has a cover and an associated heating element which form a complete thermally insulated unit. The thermally insulated unit reduces the heat loss of internally produced heat so that there is no significant loss of heat to thermally affect adjacent casseroles that may be retained in a frozen or chilled state. Such a device is not adapted to heat packaged food such as standardized “pans” size of food. Furthermore, the heating casseroles need to be regularly removed for cleaning and/or washing which may lead to potential safety issues since the casseroles also include electrical means.\nU.S. Pat. No. 5,445,062 relates to a cooker/rethermalizer especially suited for cooking or reheating of prepared, packaged meat and sauce entree items or vegetables comprising a food vessel retaining an aqueous bath, a food locator rack for supporting packaged food items and fluid outlets provided in the rack to cause fluid to exit into the bath and agitate the bath over and past food items. Such a heating apparatus is suited to accommodate flexible film packages, so called “sous vide” packages. However, the system is relatively cumbersome and requires a relatively long preheating time for the bath to be operational.\nU.S. Pat. No. 5,948,301 relates to a food thermalization device which permits the food to be rethermalized and held warm which includes an electrically-resistance heated plate which is controlled to equilibrate at a set temperature in the range of 160° F. to 185° F., with a fluctuation not exceeding plus or minus 5° F.\nU.S. Pat. No. 5,069,920 relates to a method of electric conduction cooking of food in a package. The package includes two electrodes including a raised portion of the bottom wall and a planar top electrode and wherein the recessed area between the raised portion of the bottom wall and the sidewall of the package defines a peripheral reservoir. The food is cooked by passing electrical current through the food and allowing the drainage and accumulation of exuded products from the food product being cooked so as to prevent contamination of the bottom electrode. However, the heating greatly depends upon the electrical conductivity of the food individual pieces contained in the package. As the electrical conductivity of the food product may also differ from frozen state to thawed state, and as a function of temperature, the heating pattern is difficult to control accurately. Furthermore, shock hazards are greatly increased when current is passed through the food itself whatever efforts are made to diminish the risks.\nU.S. Pat. No. 4,102,256 relates to a cooking apparatus for use in conjunction with food which is held in a container whereby, thin-walled flexible members defining the cooking surfaces are thermally coupled to heating means maintained in position against the thin-walled flexible members by maintaining means including resilient means.\nTherefore, there is a need for heating, and possibly holding hot, packaged food in relatively smaller, simpler, safer and more convenient ways using relatively low power requirements as compared to existing foodservice or even household ranges. There is also a need for efficiently and evenly heating food packaged in the standard container, in particular, aluminum lidded tray, or other differing packaging without the requirement of removing the food from the container to facilitate heating of the food to deliver the intended temperature and quality. There is also a need for controlling the heating of the food so as to optimize the heating time and reduce the risk of over heating and burnt food surfaces."}
{"text": "Embodiments of the invention relate to image processing. In particular, embodiments provide methods and systems for detecting, localizing, and classifying marker-associated points located in x-ray images."}
{"text": "1. Field of the Invention\nThe present invention relates to a hydraulic steering mechanism for steering drive wheels by utilizing hydraulic pressure.\n2. Related Art\nThere is conventionally known a hydraulic steering mechanism for steering a pair of drive wheels by connecting a pair of steering cases that are disposed at both ends of an axle case through a tie rod and by turning one of the steering cases about a kingpin shaft by a hydraulic actuator.\nHowever, in the conventional hydraulic steering mechanism, sufficient consideration has not been given to a viewpoint of a reduction in steering force of the hydraulic actuator.\nThe present invention has been accomplished with the above conventional technique in view, and it is an object of the invention to provide a hydraulic steering mechanism with a simple structure in which a hydraulic actuator for turning a steering case about a kingpin shaft can be miniaturized."}
{"text": "1. Field of the Invention\nThe present invention relates to a memory module which carries a plurality of memory elements on the surface of a laminated substrate having electrodes arranged on one side edge of the substrate.\n2. Description of the Prior Art\nA conventional memory module will be explained with reference to FIGS. 1 and 2. In FIG. 1, numeral 2 indicates a laminated substrate. A plurality of memory elements 3, for example, static random access memories (SRAMs) are mounted on both sides of the laminated substrate 2. Further, supply electrodes Vcc and grounding electrodes Gnd for the memory elements 3 are linearly arranged along the one side edge of the laminated substrate 2 together with electrodes \"In\" for signal terminals (address terminal, control terminal and data terminal). These electrodes are generally referenced as numeral 4. Such a memory module 1 is inserted into a socket provided on a circuit substrate of an electronic instrument and is thereby connected to the other electronic circuit elements making up the electronic instrument.\nWhen the removable memory module 1 is connected to its socket, electrodes for high impedance signal lines are not generally degraded in the electrical characteristics due to the configuration described above. However, electrodes for supply lines and ground lines exhibiting how impedances within the memory module 1 are often degraded in electrical characteristics, particularly, the AC characteristics arisen from the inductances and a non-linearity of the V-I characteristic accompanied by the contact with the socket.\nMore specifically, in a circuit as shown in FIG. 2, there possibly exist non-linear elements D and resistances R and inductances L caused by contamination and the like at an electrode 4 portion where the memory module 1 is brought in contact with the socket. The non-linear element D and the resistance R are not particularly active when carrying a constant current exceeding the specified magnitude. However, in the case where the supply current changes from, e.g., a non-current carrying state to a peak current carrying state, the memory elements 3 may not be supplied with adequate current. Consequently, there is a potential difference between the supply electrodes Vcc and the ground electrodes Gnd outside the memory module 1 and those inside the memory module 1.\nThe judgment as to which logic state exists should be executed on the basis of the ground level outside the memory module 1. However, this decision is actually executed on the supply and ground levels within the memory module 1. Therefore, in the case that the ground level within the memory module 1 is higher than that outside the memory module 1, a high signal input may be erroneously judged as a low signal input. Conversely, in the case that the ground level within the memory module 1 is lower, a low signal input may be erroneously judged as a high signal input.\nFor this reason, in the removable memory module, high speed memory elements have been rarely used. Also, in a memory module having a broad data bus range, when a number of bits are changed with the same timing, noise is often induced for the same reason, thus causing data error.\nA conventional memory module 1 (for example, SBX1705 type SRAM, Sony Corp.) mounting four memory elements 3 (1M bit high speed SRAM, for example, CXK 581020, Sony Corp.) on one side surface and one memory element 3 (the same as the above) on the other side surface of the laminated substrate 2 are connected to a socket provided on an electronic instrument and assembled. A schmoo plot of an address access time T.sub.AA, which gives such a noise characteristic, is shown in FIG. 3. In this figure, the ordinate shows a supply voltage V (v), and the abscissa shows an address access time T.sub.AA (ns), and also the low level input voltage V.sub.IL input in the memory elements 3 is set as 0.60 v, 0.65 v or 0.70 v. The other measurement conditions are as follows: temperature is room temperature 25.degree. C.; high level input voltage V.sub.IH input in the memory module 3 is 2.2 v; and the threshold voltage V.sub.TH is 1.5 v.\nIn this schmoo plot, the left region represents a fail region essentially pertained to the memory module 1. Another fail region caused by noise is in the center and is represented by three kinds of hatchings.\nAlso, the output enable access (hereinafter referred to as OE access) of the same memory module 1 is examined for the noise characteristic, which results in the schmoo plot as shown in FIG. 4. In this figure, the ordinate shows a supply voltage V (v), and the abscissa shows an OE access time T.sub.OE (ns), and also the low level input voltage V.sub.IL input in the memory elements 3 is set as 0.55 v, 0.60 v or 0.70 v. Even in this schmoo plot, over the whole region of T.sub.OE, there exists a fail region caused by noise at the supply voltage V exceeding 5.0 v."}
{"text": "Cutting devices have been developed which include loops for the finger and thumb of the endoscopic surgeon and an elongate stem secured to the first of the two loops at the other end of which stem is located the cutting head. A channel is usually located in the stem in which is located a wire connecting the second loop to a moveable jaw on the cutting head. The other jaw is generally fixed and formed by an aperture at the end of the stem located opposite the jaw connected to the wire. As with a pair of conventional scissors, movement of the loops away from or towards each other effects a corresponding movement in the cutting head at the end of the stem. In many surgical applications, a very small head is required for accurate cutting of body tissues. Unfortunately, there is the problem of slippage of the object being cut. This problem is caused, during surgery, because of the presence of bodily fluids reducing the friction between the cutting head and the tissue. The problem is particularly acute during endoscopic surgery because it is not possible to hold the tissue in the cutting head.\nA head has been developed with a series of lateral grooves formed on the cutting face of the moving jaw which act to hold or grip the tissue during cutting. Unfortunately, it has been found that the grooves reduce the cutting efficiency of the jaw along its cutting edges. This is particularly the case when the cutting head is very small. A cutting head has been developed with straight cutting edges along either side and a planar face which is terminated by a sharp tooth at the end of the jaw. This has resulted in improved efficiency during cutting by increasing the cutting efficiency but the problem of slippage still remains.\nAgain, this is particularly a problem where a small cutting head is required in order to cut inaccessible areas during, for instance, endoscopic surgery. In such applications, the cutting head may need to be small in order to gain access to the tissue to be cut but the tissue area may be relatively large in proportion to the size of the cutting head. Thus, it is imperative that very little slippage occurs to reduce the cutting time to a minimum."}
{"text": "Machines and machine tools are commonly provided with a series of separate, compounded slide and swivel unit arrangements to provide for multi-axis workslide linear and angular setup or positioning. Such technology corresponds to similar multiple slide/slideway arrangements commonly provided in machine tool devices for feedslide positioning as well. Particularly, very often machines are provided with a base or pedestal upon which a table is fitted for lateral movement along a predetermined axis. Surmounted on that table is generally a second slide or cross-slide provided with a second slideway arrangement for lateral movement along a second axis oriented normal to the first axis.\nAn example of common multiple slide/slideway arrangements is shown in U.S. Pat. No. 4,165,661 (which issued to A. Wasco, Jr. et al.). Wasco et al. show a machine tool comprising a bed having a pair of parallel ways fixed to the face of the bed upon which a slide is mounted for reciprocation along a horizontal path. That slide has an upper face upon which a second set of parallel ways support a second slide which provides reciprocable movement along a path normal to the first horizontal path. Other U.S. patents illustrating similar multiple slide/slideway arrangements include U.S. Pat. No. 2,627,196 (Marsilius), U.S. Pat. No. 3,347,116 (Anderson et al.), U.S. Pat. No. 3,928,943 (Wirz), U.S. Pat. No. 3,937,110 (Renoux), U.S. Pat. No. 4,159,660 (Buckley, et al.), U.S. Pat. No. 4,197,769 (Smith et al.), U.S. Pat. No. 4,719,676 (Sansone), and U.S. Pat. No. 4,813,311 (Hebbruggen).\nSimilarly, it is common to provide for multi-axis adjustment of a machine workslide by providing multiple slide/slideway arrangements. For example, U.S. Pat. No. 3,841,296 (which issued to H. Ota) describes a dressing device having a feeding table slidable on a guide surface as a result of rotative movement of a feedscrew. A traverse table is associated with and surmounts the feeding table in a slidable arrangement including a dovetail groove for movement therealong. U.S. Pat. No. 3,490,430 (which issued to N. Hoglund) illustrates an apparatus incorporating a complicated combination of a plurality of slides utilized to control the movement of a tool for dressing a grinding wheel.\nU.S. Pat. No. 4,274,388 (which issued to L. Ivel) and Pat. No. 4,071,015 (which issued to M. Funke) illustrate grinding wheel dressing devices which utilize a combination of slide/slideway structures and guide cam arrangements. Particularly, the Ivel device includes a wheel dressing assembly having a carriage moved along a linear track by a drivescrew setup, while a separate tool slide is moved in a direction transverse to that carriage by a cam and cam follower assembly. The wheel dressing device of Funke includes a cross-slide guided on dovetail ways on the dresser base, and a transversely arranged set of dovetail ways for supporting a diamond bar as it is reciprocated by the cam bar/support interaction.\nU.S. Pat. No. 4,763,549 (which issued to Y. Hata, et al.) discloses a machine tool wherein a pair of tool heads are supported on tool slides powered by feedscrews, and wherein the tool slides are moved in a second, normal direction by a servo motor attached to a second slider mechanism carried upon a separate set of ways. Such multiple-axis workslide and feedslide arrangements require additional space and machine structure to accommodate the piggy-backed or layered arrangement of slides and slideways, and require multiple joints of loading which must be independently supported.\nIn most modern day machining applications, machine stability and accuracy are critical to achieving superior products, and simplicity, compactness, reliability, and durability of the machine likewise comprise critical characteristics which must be optimized. The need for multiple load-bearing joints characteristic of the multiple tier feedslides and workslides mentioned above, adds complexity to the machine, compromising ease of manufacture, maintenance, design, reliability, stability and durability. In addition to requiring substantial redundant machine structure to support independent load-bearing slide arrangements, these independent mechanisms can further result in cumulative backlash, and more complicated design, manufacture and assembly.\nAccordingly, heretofore there has not been available in the industry a relatively simple articulated axis workslide which could provide dependable, stable, multi-directional lateral motion in a mechanism which could be easily adapted to a substantially unlimited array of machine and machine tool applications. Prior art workslide devices required a plurality of separate slides and ways to accomplish omni-directional lateral motion, or utilized compounded slide and swivel units, racks and curved gears to achieve linear and angular positioning, thereby increasing the required complexity and, generally, the size of the machine, while reducing its resulting stability, accuracy, reliability and dependability."}
{"text": "Some methods are known as the methods for producing polyhydroxycarboxylic acids of a high molecular weight by the use of Lactides; for example, a method which comprises polymerizing, in the presence of polyvalent heavy metal oxide such as lead oxide, a lactide which is obtained by heating lactic acid (U.S. Pat. No. 2,703,316), a method for producing copolymers of lactic acid and glycolic acid having an average molecular weight of 72,000 by copolymerizing, in the presence of stannous octanoate, a mixture of lactide and glycolide respectively obtained by dehydration condensation of lactic acids and glycolic acid (Japanese Published Examined Patent Application No. 14688/81), and a method for producing polylactic acids having a molecular weight of 100,000 by polymerizing DL-lactide in the presence of stannous octanoate (Japanese Published Unexamined Patent Application No. 69553/90). However, these methods, in which Lactides of high purity are necessary, require purification steps of crude lactides such as recrystallization and distillation, and thus have various problems for industrial utilization due to a number of complicated."}
{"text": "The present invention relates to a method and apparatus for converting continuous-tone images to bi-level images.\nA known method of effecting this conversion simulates the gray scale of the continuous-tone image with local variations in the proportion of black and white picture elements or pixels in the bi-level image. This is typically done by using a threshold matrix, such as a dither matrix or a density-pattern matrix, to generate blocks of bi-level pixels. A dither matrix is a matrix of threshold values that are compared one-to-one with a block of continuous-tone pixels to generate a block of bi-level pixels. A density-pattern matrix is a matrix of threshold values that are compared with a single continuous-tone pixel to generate a block of bi-level pixels. Dither matrices are used in printers, copiers, facsimile machines, and other devices to convert continuous-tone images having comparatively high resolution. Density-pattern matrices are used are used to convert continuous-tone images having lower resolution.\nUse of these methods alone, however, tends to produce a bi-level image that appears too dark. The reason is a non-linear relationship, known as a gamma curve, between perceived darkness and the actual density of black pixels in the bi-level image. Because of this relationship, the perceived darkness saturates, so that densities intended to represent different shades of dark gray look indistinguishably black.\nThe overly dark appearance and saturation can be avoided by adjusting the gray levels in the continuous-tone image according to the mathematical inverse of the gamma curve. Conventional image conversion is therefore a two-step process, comprising a gamma correction applied to the continuous-tone image, followed by the actual conversion from continuous-tone to bi-level form by dithering, density pattern generation, or a similar method.\nA disadvantage of the conventional two-step conversion process is that it takes extra time and requires extra computing resources."}
{"text": "1. Field of the Invention\nThe present invention relates to a polypropylene wrap having particular application for tampons and, more specifically, to end seals and a process for sealing the ends of a voided polypropylene film wrap.\nA tampon is enclosed in a wrap or a package prior to use to prevent the tampon from becoming soiled and also from coming into contact with moisture, either of which harmfully effects the tampon. Besides preventing the tampon from soiling and from contacting water, the wrap and its end seals need to be strong enough so as not to tear or separate during handling, yet the end seals need to tear readily to permit a user to obtain easily the tampon from within.\n2. Description of the Prior Art\nPaper has been used as a medium for a tampon wrap since it is pliable and readily permits sealing of its ends by mechanical crimping. However, a paper tampon wrap has drawbacks such as the paper can be torn relatively easily, water can easily damage the paper wrap and thus the enclosed tampon and paper does not seal out ambient moisture. Also, moisture effects paper, namely, moisture causes paper to seal. Therefore, a paper wrap is difficult to manufacture since the moisture in the wrap during manufacturing must be controlled or else the paper wrap may seal prematurely.\nThermoplastics have been used as the medium for a tampon wrap. However, thermoplastic materials, due to their nature, heretofore, could not readily be mechanically sealed and could not provide a mechanical crimped seal with the strength provided by the present invention. In particular, prior art thermoplastic tampon wraps have end seals which are formed by heat and pressure.\nSpecifically, U.S. Pat. No. 2,200,971 to Sonneborn, et al. provides a system for making, filling, and sealing packages from a thermoplastic web material in which heat and pressure are applied to form a seal. Corrugations are formed at the seal area due to the corrugations on the faces of the sealing members. This patent also provides that the end seal has beyond it a narrow unsealed margin containing a notch to facilitate opening of the package.\nU.S. Pat. No. 2,257,823 to Stokes is directed to a web which may be made of \"Pliofilm\", Cellophane, paper or other sheet material coated with, impregnated with, or consisting of thermoplastic material. The outer seal is formed by tools under conditions of temperature, pressure, and time suitable to the thermoplastic applied to or comprised in the web to thermoplastically join marginal edges of the web.\nU.S. Pat. No. 2,430,995 to Roos provides a process for sealing thermoplastic containers in which the end seals are formed by simultaneously crimping and heat sealing.\nU.S. Pat. Nos. 4,102,111 to Nack, et al. and 4,106,262 to Aterianus are directed to wrapping machines which simultaneously crimp or pressurize and heat seal a thermoplastic tube.\nLastly, U.S. Pat. No. 4,254,601 to Prager et al. is directed to an assembly for forming a plurality of individual packages each made from polypropylene. The assembly includes cooperating roller members for simultaneously heat sealing, crimping, notching, and severing a tubular plastic strip which is generally polypropylene. This patent provides that the polypropylene is heated to a temperature sufficient to melt the polypropylene so as to form a heat seal. A notch S.sup.1 is provided at each end to facilitate opening.\nWhile the prior art patents include those which use heated knurled wheels for sealing by mechanically crimping and heat sealing thermoplastics, including polypropylene, wraps and also suggest that parameters such as temperature and pressure may be adjusted as desired, see U.S. Pat. No. 2,200,971 to Sonneborn, et al., page 4, col. 1, lines 41 through 50; U.S. Pat. No. 2,257,823 to Stokes, page 2, col. 1, lines 67 to col. 2, line 2, not one of the prior art patents suggests using heated knurled wheels to form only a mechanically crimped end seal as provided by the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device having a functional circuit. In particular, the invention relates to a lightweight, thin, high-functional and low-cost semiconductor device. Further, the invention relates to a display device having a functional circuit on the same substrate, in particular to a lightweight, thin, high-functional and low-cost display device.\n2. Description of the Related Art\nIn recent years, a technology development regarding a thin film transistor (hereinafter referred to as a TFT) using polycrystalline silicon thin film as an active layer has been rapidly advanced. For example, a display device having a pixel portion of which switching element is formed of a TFT and an active matrix display device having a circuit for driving a pixel, which is formed of a TFT are actively researched and developed. Such display devices have major advantages in low-cost, thin design, lightweight, high definition, low power consumption and the like.\nFurther, research and development for forming a functional circuit by a TFT and mounting it on the same substrate as a display device are advanced as a TFT becomes higher in function. A central processing unit (CPU), an image processing unit, a memory and the like are included in a functional circuit, which are required to be formed without spoiling the advantages of a display device using a TFT as described above."}
{"text": "“Absorption” generally refers to the absorption of a wave (electromagnetic waves, sound waves, as well as light) in an absorbent material. During absorption, transmission is dampened by the material. Absorption measurement in process automation is, for example, used to determine nitrate or to measure the spectral absorption coefficient (SAC), e.g., in order to determine the organic load/charge in sewage treatment plants or in drinking water.\nFIG. 1 shows the basic measurement principle according to the prior art. The light of a light source 1, e.g., a pulsed flash lamp, irradiates the measurement section 5. Medium 5 is located in gap 6, wherein the measurement light radiated into said medium is absorbed by the determining process variable. The measurement section comprises optical windows 2, which, where applicable, also comprise lenses. A beam splitter 7 finally guides the light onto the two detectors 3 for measurement light or 3.ref for reference light. A filter 4, which allows only light of the measurement wavelength or only light of the reference wavelength to pass through in the case of the detector 3, 3.ref, is respectively mounted, where applicable, in front of the two detectors 3, 3.ref. The beam splitter 7 may also be arranged after the gap 6.\nThe above shall be explained below by way of example on the basis of a nitrate measurement. Nitrations absorb UV light in the range of approximately 190 nm to 230 nm. Nitritions have a similar absorption in the same range. In the gap 6, the nitrations and nitritions absorb UV light in the range of the measurement wavelength 214 nm in proportion to their concentration.\nIf signs of aging or temperature influences of a light source do not have any influence on the measured value of an optical probe, the light source must be monitored. In many applications, this is achieved by a second detector system (photodiode or spectrometer) directly at the light source (see FIG. 1). Any existing measurement and reference filters are in this case designed to be identical (same wavelength). If, however, this method cannot be used, either because there is no space for a second detector system in the probe or because a second detector system is too expensive, and, if a regular calibration in a zero solution also cannot be performed, the light of the light source must be guided around the measurement cuvette onto the measurement detector or the measurement spectrometer.\nFor process sensors that are to remain in the process for a long period of time, and thus also during the necessary reference measurements of the light source, the light of the light source therefore must be guided within the sensor along an optical path that is different from the measurement beam path past the measurement section onto the detector system. Since this reference beam cannot be located on the same optical axis of the measurement beam, various optical free-beam components are required in order to guide the reference beam past the measurement section and, on the receiving side, once again onto the detector system. These optical free-beam components can, for example, include lenses, beam splitters, and mirrors (see, for example, DE 100 84 057 B4). In order to generate a stable reference signal, it is also necessary to align these components correctly and robustly in relation to each other.\nThe necessary optical free-beam components occupy a lot of space, which may be an obstacle in process sensors that become smaller and smaller. Both the optical free-beam components and the production costs for their alignment are costly. Optical free-beam components can, moreover, be soiled, e.g., by humidity or by accumulation of evaporated materials."}
{"text": "Recently, for print images, the resolution in dots per inch (dpi) of printing apparatuses, such as a printer which prints digital images, is increasing along with the spread of personal computers, digital cameras, and the like. Of printing apparatuses, ink-jet printing apparatuses using an ink-jet printing head have rapidly become popular. Demands has arisen for lower-cost ink-jet printing apparatuses which realize high-resolution images.\nA serial scanning ink-jet printing apparatus prints an image of one scanning by the printing head width while scanning the printing head in a direction (main scanning direction) perpendicular to the convey direction (sub-scanning direction) of a printing medium such as a printing sheet. At the end of printing by one scanning, the ink-jet printing apparatus conveys the printing medium in the convey direction. The ink-jet printing apparatus sequentially repeats the image printing operation by one scanning described above, forming a desired image on the printing medium.\nIn a conventional printing apparatus, data sent from a host computer are data sequential in the main scanning direction (to be referred to as raster data hereinafter). To print an image of one scanning by the printing head width, data must be converted into data sequential in the sub-scanning direction (to be referred to as column data hereinafter) in accordance with the number of print elements (e.g., the number of ink discharge nozzles) of the printing head. To print an image of one scanning, the printing apparatus must hold at least a memory area for accumulating raster data by the number of discharge nozzles and a memory area for accumulating column data.\nA higher image quality is required by increasing the dpi (e.g., recording resolution) of the printing apparatus. A higher dpi of the printing apparatus increases the image resolution, resulting in a large data amount for image printing. The memory area necessary for the printing apparatus also increases, and an expensive large-capacity memory leads to high cost, inhibiting cost reduction of the printing apparatus.\nA technique of reducing the memory area necessary for the printing apparatus is disclosed in Japanese Patent Laid-Open No. 11-259248. According to Japanese Patent Laid-Open No. 11-259248, isochronous data transfer is performed between a host computer and a printer by minimum buffering. This can reduce the memory capacity necessary for the printing apparatus.\nMore specifically, image data necessary to print an image of one scanning by a width corresponding to the number of print elements of the printing head is divided into a plurality of data frames and transferred in isochronous data transfer between the host computer and the printer. When the printer receives image data transferred from the host computer, the printer processes the image data and temporarily holds it in the memory area of the printer. This memory area stores one or more data frames of image data during the printing operation by one scanning.\nIn isochronous data transfer in which image data of one scanning is divided into data frames and transferred, data frames of image data of one scanning are successively transmitted to the memory area of the printer. The memory area size of the printer suffices to ensure at least a memory area capable of storing one or more data frames, and the memory area need not store image data of one scanning. Unlike a conventional printer, a memory area for storing image data of one scanning is not required.\nIn the above-described isochronous data transfer, the printing operation by one scanning can start upon reception of minimum data necessary for image printing without storing image data of one scanning in the memory area of the printing apparatus. Sequentially transmitted data frames are stored in the memory area of the printer while the memory area is rewritten. An image of one scanning can be formed without interrupting scanning of the printing head, while the memory capacity necessary for the printing apparatus can be reduced.\nA serial printing apparatus having no memory area for one scanning of a carriage starts the printing operation upon reception of data of a given amount (e.g., one or more data frames). A memory area which stores printed data is freed, and data newly received from a host computer is stored in the freed area.\nAn example of the memory area of the serial printing apparatus and reception data in the memory area is shown in FIGS. 3A and 3B. The memory area of the serial printing apparatus comprises 0000H to FFFFH, as shown in FIG. 3A. The memory area is divided into two memory areas: a memory area (b0) of 0000H to 7FFFH and a memory area (b1) of 8000H to FFFFH, as shown in FIG. 3A. As shown in FIG. 3A, the respective memory areas are alternately used to receive and store data frames transmitted from the host computer.\nThe two memory areas are expressed as blocks (b0 and b1), b0-1th represents that the first data is received and stored in the memory area b0, and b1-1th represents that the first data is received and stored in the memory area b1. Similarly, b0-2th represents that the second data is received in the memory area b0 and the memory content is rewritten, and b1-2th represents that the second data is received in the memory area b1 and the memory content is rewritten.\nFIG. 3B shows data frames (reception data) b0-1th, b1-1th, b0-2th, b1-2th, b0-3th, . . . which are alternately accumulated in the two memory areas described above, and an image printed on a printing medium on the basis of these data frames (reception data).\nThe above printer does not have any memory area for one scanning of a carriage. The printer starts printing (print) operation upon reception of data by either of the blocks. If the other block has received data at the end of printing operation of one block, printing continues. This printing method will be explained in detail with reference to FIGS. 3C and 3D.\nFIG. 3C shows the following printing operation. More specifically, the first data (data frame) (b0-1th) is received from a host computer and stored in one memory area (b0) out of the two memory areas (blocks b0 and b1) shown in 1 of FIG. 3C (2 of FIG. 3C). Upon the completion of storage in block b0 (3 of FIG. 3C), the second data (data frame) (b1-1th) transmitted from the host computer is received and stored in block b1 as the other memory area. At the same time, image printing using the first data (b0-1th) starts (reception data decreases) (4 of FIG. 3). After the second data (b1-1th) is stored in b1 (5 of FIG. 3C), printing using the first data (b0-1th) ends (6 of FIG. 3C)\nIn addition, 7 to 12 of FIG. 3D illustrate the following printing operation subsequent to 6 of FIG. 3C. More specifically, image printing using the second data (b1-1th) starts, and the third data (data frame) (b0-2th) is received and stored in memory area b0 (7 of FIG. 3D). After the third data (b0-2th) is stored (8 of FIG. 3D), printing using the second data (b1-1th) ends (9 of FIG. 3D). Image printing using the third data (b0-2th) starts, and the fourth data (data frame) (b1-2th) is received and stored in memory area b1 (10 of FIG. 3D). After the fourth data (b1-2th) is stored (11 of FIG. 3D), printing using the third data (b0-2th) ends (12 of FIG. 3D).\nAs described above, the serial printing apparatus having no memory area for one scanning of a carriage starts printing (print) operation upon reception of data of a given amount (e.g., one or more data frames). A memory area which stores printed data is freed, and data newly received from the host is stored in the freed area. The memory area can be reduced on this premise.\nThe above-described isochronous data transfer suffers the following problem. That is, a job and task in the host computer cannot be managed by the printing apparatus. Isochronous data transfer from the host computer to the printing apparatus may be left undone due to any reason.\nFor example, if the above-mentioned data transfer from the host computer to the printing apparatus is left undone (data transfer is not in time for image printing) after the start of the printing operation by one scanning, the printing apparatus cannot print any image by the printing head. In this case, the printing head returns to, e.g., the home position (reference position) and waits until image data is stored in the memory area. After a sufficient amount of image data is stored in the memory area, the printing head is scanned again. The printing head returns from the home position (reference position) to the position where image printing is interrupted. Printing then restarts to complete image printing by one scanning.\nEvery time data transfer from the host computer is left undone, printing operation is interrupted. The printing head returns to the home position, and returns to the interrupted position in order to restart image printing, decreasing the throughput of the printing apparatus."}
{"text": "1. Field of the Invention\nThe present invention pertains to a clutch disk, comprising a hub area and a plurality of friction lining elements carried by the hub area, following each other in the circumferential direction with respect to the rotational axis of the clutch disk, these elements providing at least part of the friction surface area, where the friction lining elements are designed to make frictional contact with the associated opposing friction surface of a flywheel, a pressure plate, or the like.\n2. Description of the Related Art\nIn the clutch disks according to the state of the art, the friction linings are usually attached to a hub area, e.g., by way of a system of lining springs, in such a way that they can shift in the axial direction, as by way of the spring system, as the clutch is being engaged and disengaged. Other degrees of freedom for the movement of the friction linings are usually not available. This type of rigid connection of the friction linings to the hub area of a clutch disk, however, causes the problem that, as a result of local changes in, for example, the coefficient of friction or in the normal force prevailing between the friction linings and an opposing friction surface, frictional vibrations can be excited, which can lead to low-frequency noise, especially when relatively large masses are vibrating in phase with each other, to high-frequency noise (squeal), and even to the so-called xe2x80x9cgrabbingxe2x80x9d of the clutch.\nIt is the object of the present invention to improve a clutch disk of the general type in question in such a way that side effects produced by the frictional vibrations which occur during slipping of the clutch can be avoided as completely as possible.\nThis object is accomplished according to the invention by a clutch disk comprising a hub area and a plurality of friction lining elements supported on the hub area, which are arranged in sequence around the circumference relative to the rotational axis of the clutch disk, and which provide at least part of the friction surface area, where the friction lining elements are designed to make frictional contact with the associated opposing friction surface of a flywheel, a pressure plate, or the like.\nIt is also provided that at least one friction lining element is supported on the hub area in such a way that it can shift position in the circumferential direction relative to the rotational axis.\nIn the clutch disk according to the invention, therefore, the friction lining elements or at least one of the friction lining elements can move with an additional degree of freedom, namely, that of motion in the circumferential direction. It is therefore possible for the friction lining element to shift position slightly with respect to the hub area in response to local changes in the friction conditions between a friction lining element of this type and the opposing friction surface cooperating with it; that is, an otherwise spontaneously occurring shock is absorbed, and the danger that vibrations would be induced by such a shock is therefore almost completely eliminated.\nAs already indicated, it can be provided that the minimum of one friction lining element is designed so that it can shift position in the circumferential direction at least in response to a change in the sliding frictional interaction between its friction lining surface and the opposing friction surface which is or can be assigned to the element.\nIn order to ensure that at least one friction lining element can execute a defined movement relative to the hub area especially under the action of centrifugal force, it is also possible to provide according to the invention an additional guide device for circumferential movement, by means of which the minimum of one friction lining element is guided in the circumferential direction with respect to the hub area.\nA design which is especially reliable in operation and simple to construct can be obtained by using a telescoping guide device as the circumferential movement guide device. As an alternative, it is possible for the circumferential movement guide device, which is in either the minimum of one friction lining element or in the hub area, to comprise at least one guide recess, into which a corresponding guide projection of the other assembly, i.e., of either the hub area or of the friction lining element, can movably engage in the circumferential direction. It is also possible for the circumferential movement guide device to comprise at least one guide area or part on which a radially outward side of the minimum of one friction lining element or a support element carrying it can be supported.\nIn an alternative design of the clutch disk according to the invention, the possibility of allowing at least one friction lining element to shift in the circumferential direction relative to the rotational axis can be obtained by allowing the minimum of one friction lining element to swivel relative to the hub area around a swivel axis which is essentially parallel to the rotational axis and a certain distance away from it in the radial direction.\nTo ensure that, after a brief change in the frictional relationships, the previous set of kinematic and positional relationships are restored, it is also preferable to provide a restoring device for exerting pretension on the minimum of one friction lining element, which acts to hold it in a home circumferential position relative to the hub area.\nThis restoring device can comprise, for example, at least one spring element, which can be designed as an integral part of the hub area or of a component attached to the hub. As an alternative, it is possible for the ends of the minimum of one restoring spring element to be supported relative to the hub area and the minimum of one friction lining element.\nIn another design variant of the clutch disk according to the invention, the restoring function mentioned can be realized by providing the minimum of one friction lining element with at least two friction lining element parts, which are free to move relative to each other and relative to the hub area in the circumferential direction, and by allowing the restoring device to act between the minimum of two friction lining elements parts. In addition, by taking advantage of the dimensional stability properties of elastomeric materials, a restoring force can be obtained very easily by providing the restoring device with at least one area of elastomeric material in an area between the hub area and the minimum of one friction lining element or the friction lining of said element. In addition, it is possible for the restoring device to comprise at least one area of elastomeric material acting between two friction lining parts of the minimum of one friction lining element.\nTo take into account local changes in the frictional conditions in particular, without thus producing effects on other frictionally cooperating areas, it is proposed that it be possible for the minimum of one friction lining element or for a group of friction lining elements to move in the circumferential direction essentially independently of other friction lining elements or groups of friction lining elements.\nThe various features of novelty which characterize the invention are pointed out with particularity in the claims annexed to and forming a part of the disclosure. For a better understanding of the invention, its operating advantages, and specific objects attained by its use, reference should be had to the drawing and descriptive matter in which there are illustrated and described preferred embodiments of the invention.\nOther objects and features of the present invention will become apparent from the following detailed description considered in conjunction with the accompanying drawings. It is to be understood, however, that the drawings are designed solely for purposes of illustration and not as a definition of the limits of the invention, for which reference should be made to the appended claims. It should be further understood that the drawings are not necessarily drawn to scale and that, unless otherwise indicated, they are merely intended to conceptually illustrate the structures and procedures described herein."}
{"text": "Currently, the demand for portable devices in electronic product markets is increasing, and also, there has been a continuous demand for miniaturizing and lightening of portable devices to make them easy to carry. In order to miniaturize and lighten portable devices, packaging technology for integrating a plurality of circuit elements mounted on a printed circuit board (PCB) into a single package, as well as technology for reducing the sizes of individual electronic components provided in the portable devices are essential. In particular, a semiconductor package which deals with high frequency signals advantageously includes various electromagnetic interference (EMI) shielding structures in order to improve implementation of EMI or electromagnetic wave resistance characteristics as well as aid in miniaturization.\nTo achieve this, a related-art EMI shielding structure includes a structure which covers various circuit elements with a shield can made of press-processed metal, and a structure which forms a shielding dam made of a conductive material to enclose circuit elements, covers all the circuit elements by injecting an insulator into the shielding dam, and then forms a shielding layer thereon.\nIn the shielding structure applying the shield can, the shield can should have a constant thickness to maintain its shape, and should be spaced from each circuit element by a predetermined distance to prevent a short from the circuit element. However, due to the thickness of the shield can and the distance between the shield can and the circuit element, there is a limit to reducing the height of the shield can. Such a limit may be a factor that hinders miniaturization of the shielding structure. In addition, an air gap is formed between the shield can and the circuit element. The air gap performs an insulation action for hindering heat emitted from the circuit element from being discharged. In order to emit heat smoothly, air vents should be formed on the upper portion or side portion of the shield can. However, since electromagnetic waves leak through the air vents formed on the shield can, there is a problem that the EMI shielding effect is reduced.\nIn addition, as the technology develops, high density mounting is increasingly used. In this case, since a gap between circuit elements is set to be very narrow, it is difficult to manufacture a shielding dam satisfying a required width to height ratio by a related-art process.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid cartridge.\n2. Description of Related Art\nA known ink-jet recording apparatus is configured to record an image on a recording medium by ejecting ink stored in an ink container from nozzles. The viscosity of ink stored in the ink container may change over time. A known ink-jet recording apparatus, as described in Patent Application Publication No. JP-09-277560 A, is configured to estimate the viscosity of ink stored in an ink container, and perform optimized preliminary ejection based on the result of the estimation. More specifically, the ink-jet recording apparatus is configured to estimate the viscosity of ink based on an elapsed time since the ink container is mounted to the ink-jet recording apparatus and an amount of ink remaining in the ink container. Nevertheless, this known ink-jet recording apparatus does not estimate the viscosity by directly measuring a physical quantity obtained when ink moves in the ink container. Moreover, this known ink-jet recording apparatus cannot estimate the viscosity of ink stored in an ink container which has not been mounted to the ink-jet recording apparatus and been unused."}
{"text": "Known in the art is an apparatus for isolating titanium tetachloride from a vapour and gas mixture (cf. V.V. Sergeev, N.V. Galitsky et al. Titanium Metallurgy. Metallurg Publishers, 1971, p. 122 (in Russian)) which comprises a casing consisting of two chambers: an inlet chamber and a filtering chamber which are separated by a partition wall and provided with conical receptacles for collecting solid chlorides having chain mechanisms for chloride removal.\nThe prior art apparatus functions in the following manner.\nA vapour and gas mixture containing TiCl.sub.4, SiCl.sub.4, VOCl.sub.3, CaCl.sub.2, MgCl.sub.2, MnCl.sub.2, FeCl.sub.3, AlCl.sub.3, CrCl.sub.4, CrCl.sub.3, CO, CO.sub.2, N.sub.2, COCl.sub.2 and other vapour and gaseous chlorides is admitted at 600.degree.-950.degree. C. through a pipe to the inlet chamber where chlorides of Ca, Mn, Mg, Fe and other compounds, as well as fine particles of titanium slag are cooled, the inlet chamber and the conical receptacle being covered with a layer of chlorides and dust, the mixture flows into the filtering chamber upon an abrupt change of direction, where further cooling and precipitation of individual chlorides take place. Due to the fact that heat conductance of the deposited chlorides is poor, the throughput capacity of the apparatus decreases in time thus materially impairing quality of isolation of TiCl.sub.4 from the vapour and gas mixture.\nThe disadvantage of the prior art apparatus resides in the fact that it operates on the intermitten basis since the deposited chlorides in the conical receptacle would clogg the partition wall and the operation is interrupted as the chlorides cannot be not permanently removed from both chambers and out of the apparatus without suspending its operation. Mechanisms normally used for removing chlorides from both chambers and conical receptacles are generally inoperable since the aggressive chloride medium causes their rapid damage.\nIt is an object of the invention to provide an apparatus which eliminates heavy labour required for cleaning the conical receptacles to remove the deposited chlorides.\nAnother object of the invention is to provide an apparatus operating on the continuous basis.\nStill another object of the invention is to provide an apparatus ensuring high-grade purification of TiCl.sub.4 from other chlorides.\nThe invention consists in that in an apparatus for isolating titanium tetrachloride from a vapour and gas mixture, comprising a casing which is divided by a partition wall into two chambers; an inlet chamber and a filtering chamber for admission of a vapour and gas mixture thereto, and a conical receptacle for collecting solid chlorides arranged under each of the chambers, according to the invention, the inlet chamber and a filtering chamber are provided each with guide webs and separated by a partition wall comprising a gas duct establishing a communication between the lower portions of both chambers in such a manner that the inlet end of the gas duct is arranged under the guide webs of the inlet chamber and the outlet end is arranged under the guide webs of the filtering chamber which has inlet pipes for admission of pulps, one pipe for chloride pulp being arranged above the guide webs at a distance from one to two times the diameter of the filtering chamber and the other pipe for oxychloride pulp being arranged below the guide webs which are permanently covered by a filtering material.\nThis construction of the apparatus enables a continuous removal of the deposited chloride layer and its feeding out of the apparatus without employing heavy labour under high gas concentration, while ensuring high grade of purification of titanium tetrachloride from other chlorides.\nThe guide webs of the inlet chamber are preferably inclined in the direction of flow of a vapour and gas mixture, and the guide webs of the filtering chamber are inclined in the direction of flow of the filtering material.\nThis construction enables a better removal of chlorides from the guide webs.\nThe gas duct establishing a communication between the inlet and filtering chambers is preferably arranged in an inclined position.\nThis facility ensures a continuous flow of the deposited chlorides in a direction towards the conical receptacles.\nThe filtering material preferably comprises a carbonaceous lump sorbent, such as pitch coke.\nThis enables a better re-use of the filtering material."}
{"text": "1. Field of the Invention\nThis invention relates to certain azo spectral sensitizing dyes for photographic and photothermographic silver halide materials. In one of its aspects, it relates to photographic silver halide compositions and elements containing certain azo sensitizing dyes. In another of its aspects, this invention relates to a photothermographic element suitable for processing with heat containing photographic silver halide in reactive association with a reducing agent, a silver salt oxidizing agent, and certain azo spectral sensitizing dyes. It further relates to a photothermographic composition having the described components and a method for developing an image in such material in the absence of processing solutions.\n2. Description of the State of the Art\nVarious classes of dyes, such as cyanines, merocyanines, phthaleins, acridines, diphenylmethanes, triphenylmethanes, anthraquinones and styryls are well known as spectral sensitizers for conventional silver halide photographic materials.\nIt is also known to develop a latent image in a photothermographic element using thermal processing. After imagewise exposure, the resulting latent image in the photothermographic element is developed and, in some cases, stabilized, merely by uniformly heating the photothermographic element. Certain photothermographic materials for producing an image in color are also known. After imagewise exposure of these photothermographic materials, images can be developed by merely heating the elements to moderately elevated temperatures for a few seconds.\nIt is further known to use non-azo spectral sensitizers in photothermographic materials. These are described, for instance, in U.S. Pat. No. 3,801,321 of Evans et al, issued Apr. 2, 1974. A continuing need has existed for new classes of spectral sensitizers for photothermographic materials, especially silver halide photothermographic materials.\nAzo dyes have been useful as dyestuffs for fibers, paints, varnishes and printing inks for many years. The use of certain azo dyes in electrophotographic materials is also known.\nBritish Pat. No. 1,253,933 published Nov. 17, 1971 relates to certain azo dyes having pyrazolotriazole moieties which are useful as spectral sensitizers in conventional photographic silver halide emulsions. No use of such dyes in photothermographic materials is taught or suggested.\nU.S. Pat. No. 3,870,523 of Ikenone et al, issued Mar. 11, 1975 relates to heat developable photographic materials comprising azo dye sensitizers which have the general formulae Ar.sub.1 --N.dbd.N--Ar.sub.2 or A.sub.3 --N.dbd.N--Ar.sub.4 --N.dbd.N--Ar.sub.5 wherein Ar.sub.1, Ar.sub.2, Ar.sub.3 and Ar.sub.5 are substituted or unsubstituted aryl groups and Ar.sub.4 is arylene. This reference fails to teach or suggest the azo sensitizers useful in the photographic and photothermographic elements of the present invention.\nAzo dyes as a class are generally poor spectral sensitizers for silver halide photographic and photothermographic materials. However, they are sometimes used in photothermographic materials when conventional sensitizing dyes are unsuitable or inadequate.\nTherefore, it has been desirable to provide improved photographic and photothermographic materials comprising azo spectral sensitizing dyes, particularly when non-azo spectral sensitizers do not provide satisfactory results in such materials. It has been further desirable to use sensitizing dyes which absorb readily to the silver halide grains and which cannot be readily displaced from the grains by chemicals commonly incorporated into photographic and photothermographic materials."}
{"text": "The invention relates to a transdermal system for the administration of at least two substances through the skin with the aid of an electric current.\nTransdermal systems are generally used to administer a substance, for example a therapeutically active substance or mixture of substances, through the skin of a living organism without the need for clear penetration of the outer layer of the skin--the stratum corneum--and possibly also underlying layers of skin, by mechanical means using a device such as, for example, an injection needle. Transdermal systems are accordingly normally classed as being non-invasive dosage forms.\nThe great interest in transdermal systems stems from the fact that that dosage form has distinct advantages over other conventional dosage forms. In the case of oral administration undesired side effects frequently occur as a result of incompatibilities in the gastrointestinal tract or the liver. Also, orally administered substances are frequently decomposed or so modified in the gastrointestinal tract or the liver that the desired therapeutic effect does not occur (\"first pass\" effect). Other forms of parenteral administration, such as, for example, intravenous, subcutaneous or intramuscular injections, involve penetration of the skin or of layers of the skin and are therefore associated by the patient with a sensation of pain. In addition, local inflammation or infection may occur as a result of the partial damage to the skin. Especially in the case of long-term therapy requiring regular injections of the substance or infusions over prolonged periods--often a plurality of times a day--the patient is frequently subjected to a high level of discomfort. This has, among other things, an adverse effect on the co-operation of the patient in adhering to the medically necessary dosage scheme.\nSince transdermal systems are not subject to those limitations they are today, especially in the form of one typical example, that is to say the transdermal patch, numbered among the current dosage forms that enjoy widespread usage. Transdermal systems may be divided roughly into passive and active systems. In passive systems the substance to be administered diffuses from a reservoir through the skin. In active systems an additional force promotes the transport of substance through the skin. Especially preferred for that purpose are electric fields, which generate a flow of current through the skin. The administration of a therapeutically active substance through the skin with the aid of an electric current is generally referred to as iontophoresis.\nTypically, the iontophoretic systems used today comprise at least two electrodes, one of which forms a contact with the reservoir containing the substance. The other electrode, often referred to as the neutral electrode, is applied directly to the skin and serves to close the circuit via the body. On connection to a source of electrical energy, a current then flows through the skin and transports the substance into the body. In such active systems this is usually accompanied by a passive transport of substance.\nA disadvantage of the passive transdermal systems customary today is that the process of diffusion by way of the natural channels of the skin (sebaceous and sweat glands, inter- and trans-cellular transport paths, hair follicles) proceeds very slowly. It is accordingly difficult using a passive transdermal system to administer a substance through the stratum corneum at a dosage rate that is high enough to achieve the desired therapeutic effect. The dosage rate is the amount of substance administered through the skin per unit of time.\nCompared with passive transdermal systems, active systems, especially iontophoretic systems, usually enable higher dosage rates to be achieved. An additional advantage of the latter systems is that with them it is possible in a simple manner to influence and alter the dosage rate in a controlled manner. By regulation of the current as the active control element in the administration it is possible, for example, to adapt the dosage rate to the individual requirements of the patient. In addition, therapeutically expedient dosage schemes are feasible, for example it is possible to alternate phases of higher dosage rates with phases of low dosage rates.\nIn addition, iontophoretic systems have the advantage that the substance to be administered is, as it were, available \"on demand\". The administration of the substance can be started or stopped by simple activation or deactivation of the electric current flow.\nThere are, of course, physiologically determined limitations to the current intensities that can be used in iontophoretic systems, since too strong an electric current may result, for example, in burns or other kinds of irritation to the skin. There is therefore a need to control or increase the dosage rate by a means other than merely by way of the intensity of the electric current. One possible method of controlling or altering the dosage rate comprises administering the therapeutically active substance together with a preparation that has an effect on the blood flow in the capillaries of the skin. Such a method is disclosed, for example, in EP-A-0 448 300. In that method, first of all a composition is produced that comprises a vessel-manipulating preparation in addition to the therapeutically active substance. The joint iontophoretic administration is then carried out with that composition. If the vessel-manipulating preparation is one that dilates the vessels, then the blood flow through the capillaries of the skin is increased, resulting in a higher dosage rate of the therapeutically active substance. If the vessel-manipulating preparation is one that contracts the vessels, then the blood flow through the capillaries is reduced, resulting in a depot effect of the therapeutically active substance. Since the vessel-manipulating preparation and the therapeutically active substance are administered together, this process is referred to as \"co-iontophoresis\".\nIt is often desirable, for reasons other than for the purpose of controlling the dosage rate, to administer more than one substance iontophoretically using a transdermal system. For example different substances may have different therapeutic effects, or one substance may reduce the undesired side effects of the other substance.\nA problem with the transdermal administration of a plurality of substances using currently known iontophoretic systems, however, is that the substances, which are usually contained in the reservoir in the form of different ions, may enter into competition with one another during the iontophoresis; this means that the transport of the charge associated with the current flow is effected by means of a plurality of different charge carriers. This results in it being difficult to control the dosage rates of those substances individually. Reference is expressly made to the problem of competition also in EP-A-0 448 300 which, however, describes only the co-iontophoresis of a therapeutically active substance and a vessel-manipulating preparation. In the case of co-iontophoresis of a vessel-dilating preparation, for example, as the concentration of the vessel-dilating preparation increases, first of all an increase in the dosage rate of the therapeutically active substance is observed, but this is followed by a decrease as the concentration of the vessel-dilating preparation increases further. Consequently, it is necessary first to spend a great deal of time in determining an optimum composition of the therapeutically active substance and the vessel-manipulating preparation in order that the two substances will have the desired effect during co-iontophoresis."}
{"text": "The use of infrared radiation (IR) communications for the transmission of audio, video, data and control signals is rapidly growing. Applications using infrared transmission include remote controls for television, cable set top boxes, videocassette recorders (VCRs), digital versatile disk (DVD) players, compact disk changers and the like, remote keyboards, wireless LAN networks, video-conferencing equipment, computer peripherals, medical equipment, and personnel and equipment locating monitors.\nIn IR communications, commands and keystrokes are conveyed serially in IR packets via an IR transmission channel. The transmitted packet(s) include modulated data pulses preceded by a leader. The leader is much wider than a data pulse. The leader marks the beginning of the packet, and initiates a gain adjustment by an automatic gain control (AGC) circuit in the corresponding IR Receiver, for optimum data detection and subsequent decoding.\nBefore the rapid growth in functionality of IR remote control devices, a remote control had relatively few keys, and performance of the IR channel was not an issue. The user performed simple operations, such as: switch channel, adjust audio volume, toggle mute switch, and the like. These manual key operations were relatively slow. During a key press, a remote control device typically entered an autorepeat mode and emitted several copies of the same IR packet in a row, usually separated by an autorepeat interval. The repetition of IR Packets raised the IR channel reliability. Excess auto-repeated packets were discarded by the receiving device.\nThe appearance of more complex audio-video systems and interactive television (ITV)—in which the user utilizes an IR wireless keyboard—caused rapid saturation of the IR control channel. To meet performance requirements, typed keystrokes are now buffered in the transmitting device and are transmitted as a series of distinct IR Packets. Complex remote controls and keyboard with a multitude of keys, and pointing devices (e.g.: mouse) encode some keystrokes as a single IR packet and encode other keystrokes as a combination of several distinct IR Packets.\nSuch complex systems have been observed to suffer the problem of data loss in the same lighting conditions where simpler devices or functions still function as before. The proliferation of fluorescent lamps as a cost effective source of ambient light further degrades the reliability of the IR communication channel.\nLoss of data in the IR communication channel causes the user to repeat operations (commands, keystrokes), or choose to sit in a less desirable position much closer to the IR receiver. Errors during keyboard typing often cause marker (cursor) repositioning and necessitate retyping of lost letters on the screen. This considerably slows down typing in comparison to a (wired) computer keyboard input, drastically diminishing customer satisfaction. Some important operations are rendered difficult or impossible, e.g.: secure password entry.\nAn apparatus and method for increasing reliability without taxing performance of IR channel is desired."}
{"text": "The inventor of the present invention has discovered that it would be advantageous to provide an improved system for retaining or otherwise mounting a flexible, resilient valve defining an initially closed orifice which can be opened to accommodate fluid flow through the valve, wherein the design of the system could provide advantages not heretofore contemplated in the industry or suggested by the prior art. In particular, the system of the present invention facilitates proper closing of the valve after portions of the valve have been forced away from their initially closed configuration."}
{"text": "1. Field of the Invention\nThe invention relates to a charger and, in particular, to a two-step charger that can adjust the level of its charging command voltage depending upon the charging status of the battery.\n2. Description of Related Art\nIn view of environmental protection, many objects that cannot be decomposed by nature or cause pollution after they are decomposed are phasing out. For example, non-rechargeable batteries contain materials harmful to the environment. Therefore, rechargeable batteries become more popular nowadays. With an appropriate charger, the rechargeable batteries can be repeatedly used.\nMost of the chargers on the market have the function of two-step charging. That is, they allow the user to decide whether to charge the batteries at a normal speed or a faster speed. The so-called quick charging refers to the charging with a larger current. Although this method can finish charging in a shorter time, people normally do not immediately remove the batteries from the charger once it is done. Therefore, the larger current usually produces a high temperature on the batteries. If the batteries are exposed to the high temperature for a long time, their lifetime may be greatly shortened.\nCharging at a normal speed refers to the charging with a smaller current. As a result, the charging time is longer. It may take several hours to charge a battery in this scheme. Therefore, it is inconvenient.\nFrom the above description, it is seen that current chargers are unable to charge batteries at a quicker speed without sacrificing their lifetime."}
{"text": "This invention relates to disk drives for storing information, and more particularly, to an actuator arm assembly for a disk drive.\nDisk drives have an actuator arm which moves magnetic read/write heads into engagement with a recording media, either a flexible disk or a rigid disk in a removable cartridge.\nMost disk drives contain an actuator bearing assembly around which an arm, holding the read/write transducer, pivots. Forces (and therefore torques) are imparted to the actuator arm by means of an actuator, usually a voice coil motor, or such. These forces, besides providing the torques required to actuate the arm, tend to deflect the actuator arm and bearing assembly, and therefore, the read/write transducers, in an undesirable direction, the \"in-plane\" direction. In order to reduce such deflections, many disk drives have the top of the bearing assembly secured to a rigid structural member, such as a stiff top cover or the steel part of the magnet assembly. The bottom is usually attached to a base casting, or plate. This anchors both the top and bottom of the bearing assembly and allows for a much stiffer bearing assembly.\nHowever, as the bearing assembly is being attached to that rigid structural member, the bearing assembly can be significantly tilted so as to introduce unwanted distortion. This can be especially a problem if the stiffness or parallelism of the base casting (or plate) is poor (or competes) with the stiffness and parallelism of the other rigid structural members.\nIt is an object of the present invention to reduce the deflection of the actuator arm due to the forces applying torque to the arm, while avoiding the unwanted distortion introduced by prior art techniques of securing the bearing assembly to a structural member.\nDisk drive actuators usually have a fork-like structure at one end to attach a coil such as that of a rotary voice coil motor. Such a fork can make it difficult to locate the coil, especially in the vertical direction, while trying to attach the coil, as with an adhesive or other potting method.\nIt is another object of the present invention to provide an improved technique of locating the coil on the actuator arm."}
{"text": "The present invention relates to a package construction of a semiconductor device. More particularly, the present invention relates to a package construction of a semiconductor having a BGA (Ball Grid Array) construction where soldering balls to be used for a soldering operation in the embodying of the semiconductor device into the product are provided in a matrix on the reverse face of a substrate.\nThe conventional semiconductor device comprises a BGA substrate, a semiconductor chip to be arranged on the BGA substrate, a heat spreader for dispersing externally the heat generated in the semiconductor chip, and a ring in a space between the BGA substrate and the heat spreader for splicing both of them. The BGA substrate has a multi-layer construction with a plurality of insulation layers being superposed on each other. A plurality of lines and via holes are provided in each of the insulation layers. The BGA substrates are mutually combined with a predetermined line through the via holes when a plurality of lines are superposed. A plurality of lines can be crossed through an insulation layer, thus realizing a smaller size of the semiconductor device.\nFIG. 6 is a partially cut-away perspective illustrating view showing one example of the conventional semiconductor device. Referring to FIG. 6, reference numeral 1 denotes a BGA substrate, numeral 2 denotes a semiconductor chip, numeral 3 denotes a heat spreader, numeral 4 denotes a ring, numeral 6 denotes a soldering ball, and numeral 8 denotes a sealing member.\nEach line (not shown) provided in the BGA substrate 1 is electrically connected with the external electrode (not shown) of the semiconductor device. The soldering ball 6, made of a soldering material, is electrically connected with the external electrode of the semiconductor device. A plurality of electrodes (not shown) of the semiconductor chip 2 are electrically connected respectively with the predetermined line of the BGA substrate 1. For example, a soldering bump is provided beforehand on the surface of the external electrode connected with each electrode surface of the semiconductor chip 2 and each line of the BGA substrate 1. The connection is realized through the soldering operation by using the solder bump. The sealing member 8, made of sealing resin, is provided for adhering the semiconductor chip 2 to the BGA substrate 1. That is, sealing member 8 prevents breaking in the connecting portion between the wiring of the BGA substrate 1 and the electrode of the semiconductor chip 2 caused due to the camber or the like of the BGA substrate 1.\nIn the ring 4, an opening is provided in the center of the plate-shaped member. The shape of the opening portion is determined in accordance with the shape of the semiconductor chip 2. The shape of the heat spreader 3 is a thin plate which is similar to that of the BGA substrate 1. The semiconductor chip 2 and the heat spreader 3, the BGA substrate 1 and the ring 4, and the heat spreader 3 and the ring 4 are bonded respectively with adhesives. Adhesives for bonding the semiconductor chip 2 with the heat spreader 3 is, for example, a thermally conductive silicon adhesives. Adhesives for bonding the BGA substrate 1 with the ring 4, and the heat spreader with the ring 4 is, for example, a taper shaped thermosetting adhesive or thermoplastic adhesives.\nThen, a method of manufacturing the semiconductor device will be described. FIGS. 7(a) to 7(d) and FIGS. 8(a) to 8(c) respectively depict a process sectional illustrating view showing one example of the method of manufacturing the conventional semiconductor device. Referring to FIGS. 7(a) to 7(d) and FIGS. 8(a) to 8(c), the same reference numerals are used as the same components as those of FIG. 6. Reference numeral 5a denotes a first solder bump electrically connected with an electrode (not shown) to be included in the semiconductor chip 2. Numeral 5b denotes a second solder bump electrically connected with each one-end portion (not shown) of a plurality of lines provided in the BGA substrate. Numeral 7a shows a first adhesives layer composed of adhesives for bonding the BGA substrate 1 with the ring 4, and the heat spreader 3 with the ring 4. Numeral number 7b denotes a second bonding layer comprising a bonding agent for bonding the semiconductor chip 2 with the heat spreader 3.\nThe first soldering bump 5a is provided on the electrode included in the semiconductor chip 2. Similarly, the second soldering bump 5b is provided (see FIG. 7(a)) on each one end portion of a plurality of lines of the BGA substrate 1. A flux material is applied on the area, where the second soldering bump 5b is formed, of the surface of the BGA substrate 1. A semiconductor chip 2 is placed on the BGA substrate 1, and the BGA substrate 1 and the semiconductor chip 2 are charged into a thermally processing furnace (so-called reflow furnace) with the first soldering bump 5a and the second soldering bump 5b being in contact with each other. As a result, the first soldering bump 5a and the second soldering bump 5b are melted. Thereafter, the first soldering bump 5a and the second soldering bump 5b in contact with each other are integrated. In FIGS. 7(a) to 7(b), the first soldering bump and the second soldering bump which are integrated with each other serve as a soldering bump 5. The electrode to be included in the semiconductor chip 2 is electrically connected (see FIG. 7(b)) with a plurality of lines of the BGA substrate 1. After the flux material has been washed, the ring 4 is bonded (see FIG. 7(c)) with the BGA substrate 1 by the first adhesives layer 7a.\nAfter synthetic resin for sealing has been injected into a gap portion between the BGA substrate 1 and the semiconductor chip 2, and subjected to thermosetting to form a sealing member 8, so as to fix the semiconductor chip 2 with the semiconductor chip 2 being adhered on the BGA substrate through the sealing member 8. A second bonding layer 7b is provided (FIG. 8 (a)) by applying the adhesives on the surface of the semiconductor chip 2. The first bonding layer 7a is provided by applying the adhesives on the surface of the ring 4. Then, the heat spreader 3 is placed on the semiconductor chip 2 and the ring 4, so as to bond the heat spreader 3 on the semiconductor chip 2 and the ring 4 (see FIG. 8(b)). Finally, a soldering ball 6 is provided on the external electrode of the semiconductor device connected with the other terminal portion of a plurality of lines of the BGA substrate, so as to obtain the semiconductor device (see FIG. 8(c)).\nIn a process of manufacturing the conventional semiconductor device is a process including a thermal processing for connecting a plurality of lines of the BGA substrate with the semiconductor chip electrode by using, for example, a first soldering bump and a second soldering bump. Generally, the BGA substrate and the semiconductor chip are different in thermal expansion coefficient, so that camber is sometimes caused in the BGA substrate during the thermal processing. Therefore, the first soldering bump and the second soldering bump which have been opposed to each other before the thermal processing step are opposed to each other no more during the thermal processing process. As a result, there is caused such a problem that electric connection is not connected between a plurality of lines of the BGA substrate and the electrode of the semiconductor chip.\nOn the other hand, some unevenness exists on the BGA substrate surface and the semiconductor chip surface. When a soldering material is printed on a plurality of lines of the BGA substrate or the electrode of the semiconductor chip through the mask, so as to form the first soldering bump or the second soldering bump, the mask is not sometimes adhered with the lines or the electrodes. As a result, the size of the first soldering bump or the second soldering bump becomes unequal, thereby causing such a problem that the electric connection is not retained, between the line of the BGA substrate and the electrode of the semiconductor chip.\nAccordingly, an object of the present invention is to provide a semiconductor device which can retain the electric connection, even after the thermal processing process, between the line of the BGA substrate and the electrode of the semiconductor chip."}
{"text": "The present invention relates to a thermal ink jet printhead and method of manufacture therefor, and more particularly to an improved thermal ink jet printhead which avoids the effects of standoff between two bonded parts thereof caused by topographic formations developed during fabrication in a thick film insulating layer sandwiched between said two parts.\nIn existing thermal ink jet printing systems, an ink jet printhead expels ink droplets on demand by the selective application of a current pulse to a thermal energy generator, usually a resistor, located in capillary-filled, parallel ink channels a predetermined distance upstream from the channel nozzles or orifices. U.S. Pat. No. Re. 32,572 to Hawkins et al. exemplifies such a thermal ink jet printhead and several fabricating processes therefor. Each printhead is composed of two parts aligned and bonded together. One part is a substantially flat substrate which contains on the surface thereof a linear array of heating elements and addressing elements (heater plate), and the second part is a substrate having at least one recess anisotropically etched therein to serve as an ink supply manifold when the two parts are bonded together (channel plate). A linear array of parallel grooves are also formed in the second part, so that one end of the grooves communicate with the manifold recess and the other ends are open for use as ink droplet expelling nozzles. Many printheads can be made simultaneously by producing a plurality of sets of heating element arrays with their addressing elements on a silicon wafer and by mating a second silicon wafer having a corresponding plurality of sets of channel grooves and associated manifolds therein. After the two wafers are aligned and bonded together, the mated wafers are diced into many separate printheads.\nImprovements to such two part thermal ink jet printheads include U.S. Pat. No. 4,638,337 to Torpey et al. that discloses an improved printhead similar to that of Hawkins et al., but has each of its heating elements located in a recess (termed heater pit). The recess walls containing the heating elements prevent lateral movement of the bubbles through the nozzle and, therefore, the sudden release of vaporized ink to the atmosphere, known as blow-out, which causes ingestion of air and interrupts the printhead operation whenever this event occurs. In this patent, a thick film insulative layer such as polyimide, Riston.RTM. or Vacrel.RTM. is formed on the wafer containing the heating element and patterned to provide the recesses for the heating elements, so that the thick film layer is interposed between the two wafers when they are mated together. U.S. Pat. No. 4,774,530 to Hawkins further refines the two part printhead by disclosing an improvement over the patent to Torpey et al. In this patent, further recesses (termed bypass pits) are patterned in the thick film layer to provide a flow path for the ink from the manifold to the channels by enabling the ink to flow around the closed ends of the channels, thereby eliminating the fabrication steps required to open the groove closed ends to the manifold recess. The heater plates, having the aforementioned improvements of heater pits and bypass pits formed in the thick film insulative layer covering the heater plate surface, are aligned with and bonded to the channel plate, so that each channel groove has a recessed heating element therein and a bypass pit to provide an ink passage from the ink manifold to the channel groove.\nThorough bonding between heater and channel plates is paramount to maintaining the printing efficiency, droplet size consistency, and operational reliability of an ink jet printhead. U.S. Pat. No. 4,678,529 to Drake et al. discloses a method of bonding ink jet printhead components together by spin coating or spraying a relatively thin, uniform layer of adhesive on a flexible substrate and then manually placing the flexible substrate surface with the adhesive layer against the channel wafer surface having the etched sets of channel grooves. A uniform pressure and temperature is applied to ensure adhesive contact with all coplanar surface portions and then the flexible substrate peeled away, leaving a uniformly thin coating on the channel wafer surface to be bonded to the heater wafer. This labor intensive method tends to permit adhesive layer thickness variation between wafer pairs as well as between different regions of the same wafer pair, so that after the wafer pairs are diced into a plurality of identical printheads, the ink flow characteristics varies from printhead to printhead. A more mechanized process to place the adhesive coating on the channel wafer without operator involvement and consequent variation in parameters which introduced thickness variations in the amount of adhesive layer transferred to the channel wafers, especially in the thickness variations from wafer-to-wafer, is described in U.S. patent application Ser. No. 07/888,220, to Narang et al., Filed May 26, 1992.\nAlthough advances have improved the thickness uniformity of the adhesive layer which bonds the ink jet printhead heater and channel plates, insufficient adhesion between bonded heater and channel plates continues to cause a host of problems affecting printhead operation, such as, for example, different drop sizes between adjacent channels, because unwanted topographical protruding formations or lips are formed in the thick film layer during the patterning of the heater pits and bypass pits which prevent adequate contact between the channel wafer surface with the adhesive layer and the thick film layer on the heater wafer. Since increased adhesive layer thickness is not a practical solution because it tends to spread or wick into the channels, the inter-channel gaps between bonded heater and channel plates must be minimized in order to insure consistent printhead firing characteristics. As taught by the above identified U.S. patents, two wafers are bonded together after alignment for subsequent dicing into individual printheads. Each printhead part is formed individually on two separate substrates or wafers, where one contains heating elements and the other ink channels or passageways. The wafer containing the ink channels is silicon, and the channels are formed by an anisotropic etching process. The anisotropic or orientation dependent etching has been shown to be a high yielding process that produces very planar and highly precise channel plates. The other wafer containing the heating elements as well as heater addressing logic is covered by a thick film insulating layer in which heater and bypass pits are formed using photolithography. The thick film layer is preferably polyimide, because it is impervious to water, a major component of the printhead ink. However, one drawback with the polyimide material is its tendency to form unwanted topographical formations, such as raised edges or lips (1-3 microns high) at photoimaged edges. When bonding both heater and channel plates together, a standoff between the two plates is caused by the raised edges, which reduces the adhesiveness of the bond between the two plates and which cause the formation of inter-channel gaps.\nPolyimide topography, such as raised edges or lips, are undesirable by-products resulting from photoimaged and cured heater pits and bypass pits or trenches on heater plates. The raised edges are polyimide topographical features that are formed at the edge of photoimaged areas that do not shrink during curing as would bulk (or centered) areas of the polyimide. Consequently, raised edges critically interfere with the mating and bonding of the heater and channel plates.\nAnother form of polyimide topography is encountered in the form of edge beads or raised areas at the edge of the wafer, when a layer of liquid polyimide is dispensed and spun onto a wafer. When the contact area on the wafer is incapable of spreading further due to the contact angle at the edge of the wafer, centripetal forces push the spinning liquid polyimide towards the outside of the wafer to form an edge bead. The edge bead on a 4 inch diameter wafer, for example, is on the order of 0.5 inch wide radially from the outer edge thereof and has a thickness several micrometers thicker than the rest of the polyimide layer. Such edge beads of polyimide prevent adequate bonding between the wafers. Edge beads can also cause printhead reliability problems, because of additional stress placed on the center area of the channel plate during heater and channel plated bonding may cause cracking. Edge beads, if removed from the edge of the heater wafers, cantilevers the channel plate at its outside edges and can again cause cracks to be formed in the outer peripheral area of the channel wafer. Such cracking in the channel wafer will degrade the reliability of the individual printheads after they have separated from the wafer pair.\nRaised edges and edge beads, however, are not the only topographical formation created from photoimaged polyimide. Other topographical formations, such as wall sags or dips, compound the negative effects of raised edges by adding to the standoff between the bonded heater and channel plates. Wall dips are slumps in the polyimide walls between closely adjacent polyimide photoimaged pits. The polyimide sandwiched between the two wafers can form more than 2 microns of topographical variation, which does not allow the bonding adhesive, approximately 2 microns or less thick, to bridge or fill in the formation of inter-channel gaps. These inter-channel gaps can allow crosstalk between channels when drops are being ejected. As the patent '529 to Drake et al. teaches, care must be taken when applying adhesive in bonding the channel and heater plates so as to insure all surfaces in contact with the ink are free of adhesive, in order that the ink channels are not obstructed during operation.\nOne method of minimizing heater and channel plate standoff of printheads using a modified printhead fabrication sequence is disclosed in U.S. patent application Ser. No. 07/997,473, entitled \"Ink Jet Printhead Having Compensation For Topographical Formations Developed During Fabrication\", assigned to the same assignee as the present invention and filed on Dec. 28, 1992. The printhead enables better bonding of the two plates by compensating for raised lips or edges formed on the outside edge of opposing last pits in an array of pits located in the thick film layer that are created while photofabricating the pits in the insulating layer. The fabrication sequence compensates for the raised edges by including a non-functional straddling channel that nullifies the standoff created by the raised edge and a corresponding additional non-functional pit that positions the raised edge away from the functional channels and nozzles. Although this fabrication technique compensates for polyimide raised edges, it does not attempt to solve the problem of edge bead.\nThere continues to exist, therefore, a need to prevent the standoff between mated heater and channel plates caused by edge beads and without requiring extra non-functional, straddling channels or in drastically altering the fabrication sequence of the heater and channel plates."}
{"text": "This application relates generally to thermal ground planes. More specifically, this application relates to methods, apparatuses, and systems for flexible thermal ground planes.\nThe complexity and size of integrated circuits may been limited by the heat generated. Heat pipes have been used to transfer heat efficiently from one location to another. They have also been used to cool integrated circuits. The existing heat pipes for these purposes may consist of a rigid structure composed of copper, silicon, etc. Some modern electrical devices and systems demand a flexible circuit board along with a high capacity for heat dissipation.\nThere is thus a need for methods, systems, and devices that may also be flexible while transferring heat efficiently from one location to another or spread high flux heat from a small area to low heat flux over a larger area."}
{"text": "With the surge in the crime rate in general and, particularly in the burglary rate, there has been a proliferation of protective devices, directed mostly toward prevention of illegal entry through doors. Bars and gratings of various types have been placed across windows but, in addition to being unsightly, these constitute a serious danger in that they may prevent rapid escape in the event of fire. Moreover, the cost of such devices, when added to installation costs, makes complete protection of all windows in an apartment or home almost prohibitively expensive. Somewhat the same problem arises with respect to sliding doors whether of the type where a single door slides into a wall or the double type where either of a pair of sliding doors can be moved past the other.\nWhile complete prevention of illegal entry by intruders provides the ultimate in protection, nevertheless, an effective alarm system, which is sufficiently inexpensive so that it can be widely applied, can provide a measure of protection which substantially increases the security of the inhabitants. As is evident, then, an inexpensive but effective device, which is reliable and readily installed, is desired."}
{"text": "In dental clinics, composite resin restoration in which a caries caused in a tooth is removed to form a cavity to be thereafter filled with a resin-based filling/restoring material is widely used. The composite resin restoration is characterized by reproducing an aesthetic state similar to that of the natural teeth and allowing a treatment to be completed with only one dental visit to reduce a burden on the patient. In order to achieve aesthetic restoration, it is important to match the color tone and the transparency of the resin-based filling/restoring material with those of teeth. The teeth have a high transparency and a light color tone at incisal ends, but have a low transparency and a dark color tone at neck portions and root portions. In order to aesthetically restore, using the resin-based filling/restoring material, various portions of the teeth with different transparencies and color tones as mentioned above, it is necessary to appropriately control the color tone and the transparency of the resin-based filling/restoring material. In general, the color tone and the transparency of the resin-based filling/restoring material are controlled by adjusting the amount of an additive such as a pigment contained in the resin-based filling/restoring material. For example, Patent Document 1 proposes an adjustment method in which a composite resin is caused to express translucency using an opalescent filler.\nThe resin-based filling/restoring material itself is not adhesive to teeth. Therefore, restoration with the material inevitably requires application of a dental adhesion system. The dental adhesion system acts on the enamel and the dentin of the teeth to be interposed between the resin-based filling/restoring material and the teeth to stably bond the material and the teeth to each other. Two types of dental adhesion systems are currently available in the market: one (a total-etching dental adhesion system) involves an etching process in which a phosphoric acid is used as a pre-treatment for the teeth, and the other (a self-etching dental adhesion system) involves a self-etching process for the teeth in which a tooth substance primer containing a polymerizable monomer having an acidic group is used and no water washing is required.\nIn order to make the resin-based filling/restoration further more convenient, there are proposed various related technologies in which the operative method of the dental adhesion system is further simplified. Patent Document 2 discloses a technology (an all-in-one, self-etching dental adhesion system) in which the tooth substance can be treated with a one-step operation using an adhesive in a single package without the need for a pre-treatment performed using a pre-treatment agent such as an etching agent, a conditioner, and a primer. Further, Patent Documents 3 to 5 propose technologies in which the cavity is filled with the resin-based filling/restoring material and the resin-based filling/restoring material is cured by one step of visible light irradiation in total with no substantial visible light irradiation performed after an adhesive is applied to the surface of the cavity. However, the resin-based filling/restoring material itself is not adhesive, and a dental adhesion system is required to bond the resin-based filling/restoring material to the cavity.\nMeanwhile, dental treatments in which a glass ionomer cement that is adhesive to the tooth substance is used for filling/restoration or to seal pits and fissures are also widely performed. For the glass ionomer cement, basic glass which is a powder material is eroded by a liquid material which is an aqueous solution of an acid, and calcium ions and aluminum ions which are eluted during the erosion and carboxyl groups which are lateral groups of a polycarboxylic acid contained in the liquid material are subjected to chelate bonding and cured. Further, such ions are also subjected to chelate bonding with the tooth substance. Thus, the glass ionomer cement adheres to the tooth substance. The glass ionomer cement slightly stimulates the dental pulp, and thus can be said to be a highly biocompatible material. In addition, the material is self-adhesive to the tooth substance, and thus is widely used as a filling/restoring material (a glass ionomer-based filling/restoring material) as well as an adhesive material. The glass ionomer cement is advantageously characterized by persistently gradually releasing a minute amount of fluorine ions, and also used as a preventive material because it has a preventive effect, for example, suppressing or preventing a secondary caries and strengthening the tooth substance. In recent years, as described in Patent Document 6, there has been developed a material that has the respective characteristics of the resin-based filling/restoring material and the glass ionomer cement. Such material, however, still has the following defects: the material is low in physical properties and aesthetics compared to the resin-based filling/restoring material, and disadvantageously complicated in terms of the operative method because the material requires the concurrent use of the dental adhesion system and the pre-treatment agent."}
{"text": "The present invention relates generally to bond wire feed systems and more particularly, to an automatic feeding and threading system for a wire bonding machine.\nSemiconductor integrated circuits or dies are continually decreasing in size and therefore the diameter of the bond wires used during interconnecting semiconductor die electrodes to external leads is also shrinking. As a result of this shrinking, manual wire threading during bond wire spool changing, or wire bond trouble shooting, becomes more time consuming for operators.\nIn addition to the reduced diameters of bond wires, there is a move towards palladium coated copper bond wires, which are difficult to distinguish from the background of a bonding machine. As a result, damage to die electrodes may occur when wire bonding is attempted when an undetected insufficient bond wire length protrudes through the capillary welding tube. This insufficient bond wire length (short tail or no tail) is caused when the bond wire does not protrude sufficiently from the end of the capillary bonding head after a successful bond.\nTherefore, it is an object of the present invention to at least alleviate one of the problems associated with bond wire feeding."}
{"text": "Chemical messengers are small, diffusible molecules within living organisms and include second messengers such as ions, hormones, neurotransmitters, cyclic nucleotides etc that play a central role in development and cell function. Methods that are generalized to capture maps of chemical messengers within cells are invaluable in order to understand the manifold functions of second messengers. Methods to obtain chemical maps specific to a given messenger mostly use genetically encodable fluorescent sensors.\nIn contrast, methods that are generalizable to multiple second messengers tend to use imaging based on CARS, SIMS and MALDI, still need to achieve either adequate temporal resolution or reduce imaging artifacts. Thus obtaining high resolution spatiotemporal chemical maps in living systems using a generalizable methodology is still challenging. A second messenger that plays a crucial role in metabolism, neuronal activity, cell-cycle control and growth is pH. Different second messengers are functionally coupled and act in concert to stringently maintain sub-cellular proton concentrations since pH regulates the activity of key enzymes and ion channels.\nSubcellular organelles are bounded by membranes where the intra-organellar pH is stringently regulated. The lumenal pH of various cellular organelles has been mapped primarily using pH-sensitive fluorophore functionalised ligands that bind specifically to receptors that are resident in the relevant organelles. As the receptor shuttles between the plasma membrane and its relevant organelles, the derivatised ligand is ferried by the receptor along its retrograde endocytic pathway. Organelle pH may also be measured by expressing pH sensitive fluorescent proteins fused to peptide sequences that function as organelle localization signals. However these cannot give temporal information on pH changes of a receptor containing compartment while it undergoes maturation.\nThus, most studies have utilized pH reporters in the form of fluorescent pH probes conjugated to endocytosable ligands. However, for the vast majority of cases where (a) the fluorophore conjugation to the ligand either disrupts ligand structure or ligand trafficking and (b) for membrane proteins which have no associated ligands, pH mapping of the relevant compartmental maturation process is still not accessible.\nThe burgeoning field of DNA nanotechnology has yielded a number of powerful synthetic molecular devices for small molecule sensing in vitro. Yet remarkably, this chemical diversity in sensing has not yet been exploited in cellulo or in vivo. One of the rare examples of DNA-based molecular devices that show quantitative preservation of its sensing functionality both in cellulo and in vivo, is the I-switch. This is a DNA assembly that undergoes a conformational change triggered by acidic pH. Acidic pH causes the formation of a non-Watson-Crick based DNA motif called the I-tetraplex, or i-motif that is then transduced into a large scale conformational change of the overall DNA assembly. The I-switch has therefore been used in living systems as a pH sensor to map spatial and temporal pH changes associated with the maturation of endosomes, by conjugating it with endocytic ligands. However, such DNA devices are not amenable to report on the chemical environments of the vast majority of proteins.\nDNA has been used to build nanomechanical devices with potential in cellulo and in vivo applications. However, their in cellulo applications in different biological pathways are limited due to current device response times as well as limitations associated with their delivery to precise intracellular locations."}
{"text": "There is known a NAND flash memory in which memory cells are three-dimensionally arranged."}
{"text": "Virtual personal assistants are artificial intelligence systems that perform tasks on a computing device in response to natural-language requests from a user. Typical virtual personal assistants are monolithic applications that handle a finite set of tasks for the user, such as calendaring, reminders, and messaging. Those virtual personal assistants understand how to respond to a number of natural language phrases associated with those tasks. Extending the capabilities of virtual personal assistants generally requires extensive modifications to analyze and respond to natural language relating to the new capabilities. For example, systems based on a grammar model must be substantially changed for any new vocabulary or phrasing behavior. Similarly, systems based on a dictation model may require extensive effort to integrate required training sets of spoken data. The effort and/or amount of data required to integrate new capabilities may increase exponentially as additional capabilities are added."}
{"text": "Various colloidal particles of alumina hydrates have already been produced by various processes for production. The colloidal particles are a generic term for particles having a particle size of 1 to 1000 nm, that is, those having a particle size of colloidal level. Most of these produced colloidal particles are composed of alumina hydrates having a boehmite structure or pseudoboehmite structure, or amorphous alumina hydrates, and it is known that they are plate-like, ribbon-like, spindle-like, needle-like or fibrous in shape.\nJP 57-111237 A (1982) discloses a process for producing an alumina sol comprising hydrothermally treating an alumina hydrate obtained by reacting a water-soluble aluminum salt with carbon dioxide or a carbonate to obtain a processed product, and then drying it, or drying the processed product, and then mixing it with a monobasic acid.\nJP 02-243512 A (1990) discloses a process for isolating a fine and plate-like α-aluminum monohydrate (alumina hydrate having a boehmite structure) product by controlling molar ratio of a caustic aluminate liquor at a temperature below 100° C. and cooling the alkaline solution with vigorous stirring.\nU.S. Pat. No. 4,666,614 discloses a process for producing α-aluminum monohydrate comprising reacting an acid or alkali reaction solution controlled to pH of 7.5 to 10 at 60 to 100° C. for 2 to 7 hours and drying it to 35 mass % or more.\nJP 57-111237 A (1982) states that as to production condition of alumina hydrate, in case where an aluminate of alkali metal is used as a water-soluble aluminum salt, it is sufficient to follow a usual production process in which gaseous carbon dioxide is bubbled into the salt and the carbon dioxide is desirably introduced until the reaction system reaches pH near 7, and then the produced alumina hydrate is washed for removing impurities. Further, the publication states that smaller amount of remaining impurities is preferable from viewpoint of the production or use of alumina sol.\nJP 02-243512 A (1990) states that the α-aluminum monohydrate of the invention is plate-like crystal having a size (average particle size is 0.4 μm, or in particle size distribution, 90% or more of particles are in a range from 0.2 to 0.8 μm) suitable for pigment filler, and that boehmite separating maximum rate from sodium aluminate solution is obtained at a molar ratio Na2O/Al2O3 of 1.3.\nU.S. Pat. No. 4,666,614 states that an aluminum compound is reacted with an acid or an alkali, pH is adjusted to 7.5 to 10, and then heating is carried out at 60 to 100° C., thereafter by-products formed by neutralization are removed to obtain a sol having an alumina hydrate concentration of 10 to 35 mass %.\nIn the production processes disclosed in JP 57-111237 A (1982) and JP 02-243512 A (1990), the need of washing step carried out after the neutralization with an alkali and an acid is stated. However, the washing is carried out with pure water, not only it is time-consuming until completion but also the intermediates allowed to stand after washing and filtration tend to cause aluminum hydroxide. The aluminum hydroxide is produced due to aging of the aimed product (amorphous alumina gel). And, alumina hydrate having a boehmite structure originally aimed is not obtained at a normal temperature. In order to prevent the aging, carbonate or gaseous carbon dioxide, etc. is generally added. However, as carbon dioxide is weak acid, a large amount of it is required for neutralization, therefore it is economically disadvantageous to adjust to pH around neutrality.\nIn addition, in case where the alumina obtained by the process disclosed in JP 02-243512 A (1990) is used as a filler for filling pigment, it cannot show its sufficient performance for any application purpose as the size of the particles is large. Therefore, those having further small particle size are required.\nIn the production processes disclosed in JP 57-111237A (1982) and U.S. Pat. No. 4,666,614, a large amount of acid or alkali used at pH range in neutralization step is required and thereby leading to high production cost, and further there are some problems that a large amount of pure water or chemicals and time are required for washing by-products, and efficiency is not good.\nFurther, the required physical properties such as viscosity and thixotropic property of the aqueous alumina sol are different according to each application purpose. It is desired to provide a further improved aqueous alumina sol for specific application purposes."}
{"text": "At present, banknote boxes in financial self-service equipment (such as ATM) can be classified into stack-type banknote boxes and drum-type banknote boxes according to different ways in which banknotes enter and leave the banknote boxes. Drum-type banknote boxes have been applied in circulatory machine cores due to its advantages such as having a small occupation volume, a high reliability of banknotes entering and leaving, and a considerable amount of banknotes deposit and so on. The depositing of banknotes in the banknote box is completed by a reel and a tape, and the banknotes are guided by the tape to wrap around the reel layer upon layer. When the banknotes are discharged, the tape rotates reversely to guide the banknotes to be released from the reel.\nWhen the financial self-service equipment is used less frequently, the banknotes may appear to have a curvature approximate to a diameter of the reel after having been stored in the reel of a drum-type machine core, and especially the banknotes near an inner core of the reel have the largest curvature. Therefore, when discharging and conveying severely curled banknotes, unsmooth conveying and even jamming and tearing of the curled banknotes may easily occur at a turning and direction changing position in a passage; besides, curled banknotes are not conducive to a secondary banknote feeding, because ends of the banknotes are apt to be bent which obstructs the banknotes from entering a banknote entrance, and thus severely affects the reliability of equipment operation and customer experience.\nAt present, a common method to flatten a banknote is to make the banknote pass through a passage having a bending direction opposite to a direction in which the banknote is bent, to force the banknote to bend and deform reversely, and alleviate a bending and deformation degree of the banknote through time and pressure. However, due to a high transmission speed of banknotes in the financial self-service equipment, the treated banknotes have an unobvious or even no flattening effect; and the efficiency is low, which fails to meet the application requirements of the financial self-service equipment."}
{"text": "Aspects of this disclosure relate generally to telecommunications, and more particularly to techniques for sounding sequence protection in multi-user multiple-input-multiple-output (MU-MIMO) communications for wireless local area networks (WLANs).\nThe deployment of WLANs in the home, the office, and various public facilities is commonplace today. Such networks typically employ a wireless access point (AP) that connects a number of wireless stations (STAs) in a specific locality (e.g., home, office, public facility, etc.) to another network, such as the Internet or the like. A set of STAs can communicate with each other through a common AP in what is referred to as a basic service set (BSS). Nearby BSSs may have overlapping coverage areas and such BSSs may be referred to as overlapping BSSs or OBSSs. In some scenarios, communications that occur in nearby BSSs can result in collisions and failure in the transmission of information.\nIn dense enterprise deployments of WLANs, such as in stadiums, airports, or other large venues, for example, there may be multiple APs deployed, and the coverage of several of those APs can overlap creating OBSS scenarios. For example, in these dense deployments, multiple STAs can be in the common coverage of multiple BSSs. Moreover, when these dense deployments are unplanned, some of the APs may be automatically configured to work on the same channel, which may cause transmission collisions between OBSSs. The collisions that occur may result in sounding sequence failures and, upon detection of a sounding sequence failure, an AP may terminate a transmission opportunity (TxOP) and would need to contend for the medium again. System throughput can be severely impacted if this happens frequently.\nFor MU-MIMO transmissions, however, the WLAN standards (e.g., IEEE 802.11-based standards) have not defined a specific mechanism that may be used for MU-MIMO sounding sequence protection. One option for MU PLCP Protocol Data Unit (PPDU) protection is to send multiple RTSs one by one to all of the STAs that will be part of the subsequent MU PPDU transmission and expect a CTS from each of those STAs. However, this solution may not be practical for MU-MIMO sounding sequence protection because it incurs a large overhead that may not be needed in non-OBSS scenarios.\nAccordingly, for scenarios that create OBSSs between nearby BSSs and that can result in sounding sequence collisions and failures, it may be desirable to have a mechanism that protects the sounding sequence while providing low overhead."}
{"text": "1. Field of the Invention\nThe present invention relates to a detachment-preventing plug; in particular, the invention relates to a detachment-preventing plug without a corresponding socket.\n2. Description of the Related Art\nConventional detachment-preventing devices are often disposed on sockets rather than on the plugs. For example, FIG. 1 shows a conventional locking frame disposed outside the socket, and the frame defines a compartment for receiving a plug. The drawback of such conventional design is that the compartment defined by the frame must correspond with the size of the plug, or the frame can't strengthen the fixing capability. Furthermore, the frame occupies a certain space, and the usages of the electric equipments are restricted within the place having the special socket. In other words, the electric equipments do not have the detachment-preventing function, which are carried by a user to the place without the special socket.\nBased on research, the inventor proposes the present invention to address the above issues."}
{"text": "1. Field of the Invention\nThe subject invention is directed to gas turbines, and more particularly, to a system for delivering fuel to the combustion chamber of a gas turbine engine by lean direct injection.\n2. Background of the Related Art\nWith increased regulation of pollutants from gas turbine engines, a number of concepts have been developed to reduce engine emissions while improving engine efficiency and overall operability. One such concept is the use of staged combustion. Here, the combustion process is divided into two or more stages or zones, which are generally separated from each other, either radially or axially, but still permitted some measure of interaction. For example, the combustion process may be divided into a pilot combustion stage and a main combustion stage. Each stage is designed to provide a certain range of operability, while maintaining control over the levels of pollutant formation. For low power operation, only the pilot stage is active. For higher power conditions, both the pilot and main stages may be active. In this way, proper fuel-to-air ratios can be controlled for efficient combustion, reduced emissions, and good stability.\nIn addition to staged combustion, providing a thoroughly blended fuel-air mixture prior to combustion, wherein the fuel-to-air ratio is below the stoichiometric level so that combustion occurs at lean conditions, can significantly reduce engine emissions. Lean burning results in lower flame temperatures than would occur during stoichiometric combustion. Since the production of NOx is a strong function of temperature, a reduced flame temperature results in lower levels of NOx. The concept of directly injecting liquid fuel into the combustion chamber of a gas turbine and enabling rapid mixing with air at lean fuel-to-air ratios is called lean direct injection (LDI).\nThe prior art is replete with example of LDI systems. For example, U.S. Pat. No. 6,389,815 Hura et al. discloses a lean direct injection system, which utilizes radially staged combustion within a single injector. The pilot fuel delivery stage includes a pressure swirl atomizer that sprays liquid fuel onto a filming surface. The liquid film is then stripped off into droplets by the action of compressor discharge air. The main fuel delivery system includes a series of discrete atomizers that spray fuel radially outward into a swirling cross-flow of air. The main fuel delivery system is staged radially outboard of the pilot fuel delivery system, and operates in the fuel-lean mode. Radial separation as well as an air jet located radially between the two stages achieves separation of the pilot combustion zone and the main combustion zone.\nU.S. Pat. No. 6,272,840 Crocker et al. discloses a lean direct injection system, which also utilizes radially staged combustion within a single injector. The pilot fuel delivery is either a simplex air-blast type atomizer or a prefilming air-blast type atomizer, and the main fuel delivery system is a prefilming air-blast type atomizer. Separation of the pilot and main combustion zones is achieved by providing an air splitter between the pilot outer air swirler and the main inner air swirler. The air splitter develops a bifurcated recirculation zone that separates the axially aft flow of the pilot injector from the axially aft flow of the main injector. The bifurcated recirculation zone aerodynamically isolates the pilot flame from the main flame, and ensures that the pilot combustion zone remains on-axis with no central recirculation zone. A converging wall of the pilot air cap, which essentially acts as a flame holder to anchor the flame, defines the air splitter. Acting in this manner, the pilot air cap will likely suffer thermal distress (i.e., oxidation, melting), and require some form of thermal management. In this regard, Crocker et al. disclose the use of small cooling holes in the air cap to improve durability.\nEuropean Patent Application EP 1413830 A2 discloses a lean direct injection system, which also utilizes radially staged combustion. In this case, an air splitter with an aft end cone angled radially outward assists in creating a bifurcated recirculation zone. The additional function of the splitter is to prevent the inner main air stream from modulating with combustor pressure fluctuations, thus reducing combustion instability. This air splitter has a larger radial extent than the air splitter disclosed in U.S. Pat. No. 6,272,840 to Crocker et al., and acts as an even larger flame-holder, requiring thermal management to avoid distress.\nWhile the concept of the LDI system is sound, achieving the required levels of performance can be difficult. Lean-burning systems are prone to localized flame extinction and re-ignition. This results in combustion instability that can damage the combustion chamber. Limitations in atomization, vaporization, and fuel-air mixing can result in heterogeneous stoichiometric burning, which yield higher than desired levels of NOx. Also, for these self-contained radially staged LDI systems, control over the level of mixing between the pilot combustion zone and the main combustion zone can be difficult. The negative effects can include reduced margin for lean blowout, and possibly increased levels of smoke.\nAccordingly, there is a continuing need in the art to provide a lean direct injection system which can achieve low levels of combustion instability, enhanced atomization quality, increased fuel-air mixing rates, low pollutant formation, low smoke and improved lean blow-out margin."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent the work is described in this background section, as well as aspects of the description which may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nReferring now to FIG. 1, conventional cable testers 10 are frequently used to isolate cabling problems and to measure cable-lengths. The cable tester 10 is coupled to a cable 14 by a connector 12. The connector 12 may include an RJ-45 or other suitable connector. A connector 15 connects the cable to a load 16. The cable tester 10 typically uses a loop-back module (not shown) as the load 16 at a remote end. The cable tester 10 performs cable analysis and can detect faults—e.g., a short, an open, a crossed pair, or a reversed pair—in the cable 14. A short or an open in the cable 14 may be detected without a load. The cable tester 10 can also determine a length of the cable 14 and a distance from one end of the cable 14 to a point where the cable 14 has a fault such as a short or an open condition.\nFor example, in a multi-conductor cable, a short condition occurs when two or more conductors in the cable 14 are short-circuited together. An open condition occurs when one or more conductors in the cable 14 lack continuity between both ends of the cable 14. A crossed pair occurs when a pair of conductors communicates with different pins at each end of the cable 14. For example, a first pair of conductors may communicate with pins 1 and 2 at one end and pins 3 and 6 at the other end. A reversed pair occurs when two ends in a pair are connected to opposite pins at each end of the cable 14. For example, a conductor connected to pin 1 on one end communicates with pin 2 at the other end, and a conductor connected to pin 2 on one end communicates with pin 1 at the other end.\nThe cable tester 10 typically employs time domain reflectometry (TDR) (which is based on transmission line theory) to troubleshoot cable faults and to measure cable-lengths. In operation, the cable tester 10 transmits a test pulse 17 on the cable 14 and analyzes a corresponding reflection or a return pulse 18. Specifically, the cable tester 10 measures a difference between a time when the test pulse 17 is transmitted and a time when the return pulse 18 is received. Additionally, the cable tester 10 analyzes characteristics such as shape and size of the return pulse 18 relative to the test pulse 17. Thus, a fault in the cable 14 as well as a length of the cable 14 can be determined based on electrical properties of the cable 14 (e.g., a cable propagation constant) and the comparisons above between the test pulse 17 and the return pulse 18.\nA conventional cable tester (e.g., cable tester 10), however, may generate inaccurate results when a cable is properly terminated at the remote end. For example, TDR techniques typically cannot be used to determine a length of cable when the cable is connected (at the remote end) to a link partner that is active or in use. In such a case, the cable functions as a substantially balanced transmission line when the remote end of the cable is properly terminated. With a substantially balanced transmission line, when the remote end receives a TDR pulse (e.g., test pulse 17), the remote end may return a very weak signal. Weak return signals typically cannot be analyzed unless extensive electronic circuits are used. If the cable is perfectly terminated, there may be no reflected signal. Implementing extensive electronic circuits, however, can be expensive and may not be feasible in low-cost systems.\nAlternatively, digital signal processing (DSP) techniques can be used to determine a length of a cable when the cable is connected to an active link. In DSP, unlike in TDR, no pulses are injected into a cable. Instead, parameters such as amplitude, pulse width, pulse shape, etc., of signals that are normally transmitted and received on the cable are measured to determine a length of the cable (or cable-length). DSP, however, involves making some assumptions and therefore yields cable-length measurements that are approximate rather than accurate.\nFor example, if the length of a cable is determined based on an amplitude of a received signal, the amplitude of the transmitted signal is generally unknown or unknowable and, therefore, needs to be assumed. Additionally, any attenuation in the received signal is calculated by assuming an average attenuation per unit length of the cable. Therefore, the length of a cable determined using DSP techniques is generally an approximate estimate rather than an accurate measurement.\nThus, in low-cost systems, since TDR techniques generally cannot be used to analyze reflections from a properly terminated end or an active link, TDR techniques cannot be used to determine a cable-length although the cable may have no faults. Furthermore, while DSP techniques can be used to determine a length of a cable that is properly terminated or that is connected to an active link, DSP techniques generally cannot be used to determine an accurate cable-length when the cable is very long."}
{"text": "The present invention relates to carton erecting machines, and, more particularly, to such machines which employ vacuum devices for gripping both surfaces of the folded carton and a mechanism for moving one of the vacuum devices with respect to the other to pull the carton into an opened condition.\nFolded cartons are formed from flat carton blanks having score lines defining the individual panels of the carton and having a narrow tapered strip or edge on one end for glueing. The blank is bent at each of score lines defining the side edges of the panels, and the tapered strip is glued to the edge of the panel on the opposite end of the blank. The folded carton, when properly formed, is a tubular parallelogram which flattens out under the weight of a stack of cartons but tends to spring open again when pressure is removed. Such cartons are easy to open, and if made of board of uniform thickness and stiffness, can be erected by the conventional prior art machines without incident.\nIn actual practice, a significant percentage of folded cartons are not properly and uniformly made and the prior art machines are not capable of effectively erecting such imperfect cartons. The machine therefore malfunctions causing a temporary shut down of the line while the machine is cleared of faulty cartons or readjusted to accomodate cartons made of different board.\nOne common fault in carton construction is failure to bend the blank along all appropriate score lines as it is formed. This makes the carton difficult to open since the opening mechanism must supply enough force to bend the carton along that score line. A second common fault is the partial glueing together of the internal surfaces of the carton caused by misalignment of the glueing mechanism or the application of an excessive amount of glue. In order to open these cartons, this glue bond must be broken by the erecting mechanism. A third common problem is warpage of the cartons due to improper storage or construction. Warpage has two adverse effects. When it bends the carton transversely to the score lines, it makes opening of the carton difficult, and it changes the position of the surface of the end carton in the magazine causing problems in removing the carton from the magazine.\nAnother factor which causes many prior art carton erecting machines to malfunction is a variation in the stiffness of the cartons due to a change in the board from which they are formed. Where the carton at the open end of a magazine is held in by lips which engage opposite edges of the carton, and where the magazine is oriented so that the weight of the cartons act against the end carton, the end carton bends in an outward arc between the two retaining lips. The extent of this deflection is dependent upon the stiffness of the board from which the cartons are made. The prior art erecting machines using this type of magazine must be adjusted according to the extent of the deflection and therefore a variation in the stiffness of the cartons causes the machine to malfunction and requires readjustment."}
{"text": "This invention relates to therapeutic peptides.\nA number of luteinizing hormone releasing hormone (LH-RH) analogs have been described which inhibit the release of LH-RH, a peptide hormone having the formula pGlu-His-Trp-Ser-Tyr-Gly-Leu-Arg-Pro-Gly-NH.sub.2. For example, Coy et al. U.S. Pat. No. 4,431,635, hereby incorporated by reference, describes LH-RH analogs having the general formula X-R.sup.1 -R.sup.2 -R.sup.3 -Ser-Tyr-R.sup.4 -Leu-Arg-Pro-R.sup.5 -NH.sub.2, in which X can be Acetyl (AC); R.sup.1 and R.sup.4, independently, can be D-Trp or D-p-X-Phe, where X is a halogen or methyl group and (represents such throughout the rest of this application); R.sup.2 can be D-p-X-Phe; R.sup.3 can be D-Trp; and R.sup.10 can be Gly or D-Ala."}
{"text": "Detector devices for receiving light and generating electrical signals, including a housing and a detector, often exhibit a noise-inducing, temperature-dependent dark current. Such a dark current can be reduced by cooling.\nThe German Patent Application DE 10 2009 036 066 A1 discusses an optoelectronic detector that has a cooling device, namely a Peltier element, that is thermoconductively connected to the detector. To prevent condensation water from forming on a surface of the optoelectronic detector, a sensor is provided for determining an instantaneous value of the ambient air humidity and of the ambient dew point temperature. The sensor is connected to a control unit that controls the cooling device as a function of the value. It is advantageous that this optoelectronic detector does not dispense altogether with a cooling. However, the actual cooling capacity is disadvantageously limited to low levels, namely to levels at which no condensation water forms. Detector noises are, therefore, not effectively prevented.\nThe same publication mentions another detector device where the detector, together with the cooling device, typically a Peltier element, is encapsulated in an airtight housing that is filled with a dry gas or is evacuated. In the case of this device, the waste heat can be transferred to a heat sink that is thermoconductively connected to the cooling device, and/or be used for heating other components, such as an entrance aperture of the housing. However, this detector device is recognized as being disadvantageous because of the costly air-tight encapsulation.\nIn practice, it actually turns out that other disadvantages are associated with this detector device. In particular, the cooling is often not very effective. Moreover, the cooling proves to be especially problematic when the detector is required to be at a different electric potential level than the housing. In such a case, the Peltier element cannot be readily configured between the housing and the detector. Such a potential difference is mostly necessary when photoelectrons are to be accelerated within the detector."}
{"text": "1. Field of the Invention\nThis invention relates generally to drilling mud separation systems, and more particularly to a drilling mud separation system which utilizes a moving continuous belt in the form of a screen, chain, self-cleaning chain link belt, or combination chain link belt and wire mesh screen to remove large drill solids and gumbo from the drilling mud or drilling fluid prior to the drilling fluid or drilling mud flowing to conventional solids control equipment.\n2. Brief Description of the Prior Art\nEffective solids control systems for controlling the solids in drilling mud has been a goal in the oilfield industry for many years. Various have been employed to give higher penetration rates, less contamination of formations, lower mud costs, and less abrasives and sticky materials in the mud system. The systems which are most effective require expensive and complicated machinery and require personnel to monitor and maintain the machinery. Common apparatus used for mechanically removing solids from drilling mud include shale shakers and vibratory screens, desanders, desilters, mud cleaners and centrifuges. Each piece of equipment is limited to a range of various particle sizes and selectively reject undesirable solids and retain desirable solids in the drilling fluid.\nThe shale shakers and mud cleaners employ a vibrating screen to selectively classify particles by size differences. The desander, desilter and centrifuge are usually located downstream from the shale shakers and utilize centrifugal force and mass difference between the solids density and liquid density for solids removal.\nThe initial removal step in most solids control systems is to conduct the drilling mud from the well directly to a shale shaker. Vibrating screen shale shakers of the type used in the oil industry fall into three general categories, as described below.\n\"Elliptical motion\", or \"unbalanced design\" shakers have a downward slope on the screen to transport cuttings across the screen and off the discharge end. Optimum screening with these types of shakers is usually in the range of 30-40 mesh (400-600 microns), and they are prone to bearing failure.\n\"Circular motion\", or \"balanced design\" shakers produce a balanced, or generally circular motion. The consistent, circular vibration allows solids transport with the basket in a flat, horizontal orientation. This type of shaker usually has multiple decks to split the solids load and allow the use of finer mesh screens in the range of 80-100 mesh (150-180 microns).\n\"Linear motion\", or \"straight-line motion\" shakers produce a generally straight motion. This motion is developed by a pair of eccentric shafts rotating in opposite directions. Linear motion shakers provide superior cutting conveyance and are able to operate at uphill slope which allows the use of 200 mesh screens (77 microns).\nMost prior art shakers tend to force the cuttings upward from the screen surface and as they continue upward the screen travels down, and as the cuttings fall downward, the screen travels upward and strikes the cuttings which chips small pieces off and creates very fine cuttings that are extremely difficult to remove. Shale shakers are prone to short screen life and their vibratory action tends to drive the cuttings into the screen and leads to screen \"blinding\" resulting in loss of fluid across the screen. \"Blinding\" is known as the phenomenon where near-size particles plug the screen or the screen becomes coated with sticky particles.\nThe present invention is not a replacement for the conventional solids control equipment but instead is installed upstream from the conventional solids control equipment and removes large amounts of large drill solids and gumbo from the drilling mud or drilling fluid prior to the drilling fluid or drill mud flowing to the conventional solids control equipment such as shale shakers. The present apparatus utilizes a separation unit removably installed in a box-like housing which receives the drilling mud or drilling fluid from the well separates the large drill solids and gumbo by utilizing a moving continuous belt in the form of a screen, chain, self-cleaning chain link belt, or combination chain link belt and wire mesh screen which moves in an continuous loop. By removing large drill solids and gumbo from the drilling mud or drilling fluid, the present apparatus improves the efficiency and performance and reduces the cost of operating and maintaining the conventional solids control equipment. Removably mounting the separation unit in the box-like housing allows a replacement separation unit to be easily and quickly installed so that operations may continue while the removed separation unit is repaired."}
{"text": "Not Applicable\nThe present invention relates to pipe systems and more particularly to furnace pipe systems and more particularly to liquid evacuation devices for furnace pipe systems.\nDrainage of water condensation from furnace exhaust pipes in typical heating systems accumulates in undesirable areas and causes deterioration or damage to heating system components. Water condenses from the atmosphere onto the exhaust pipe walls because the hot gasses which travel in the exhaust pipe contact the exhaust pipe walls which are often much cooler. The difference between the temperature of the hot gasses and the temperature of the exhaust pipe walls causes water condensation to form on the exhaust pipe walls because the quantity of water vapor which may be contained in a volume gas decreases as the gas temperature decreases. When the temperature of a gas is decreased by contact with a cooler surface the gas may become saturated with water vapor so that water will condense out of the gas and onto the cool surface as liquid. Water condensation typically drains along the exhaust pipe walls and accumulates in pools on any horizontal surface. It is desirable to prevent such accumulation of water and thereby extend the useful live of heating system components.\nU.S. Pat. No. 5,620,302 teaches that drainage holes which are provided in a furnace blower housing allow water to drain out of a furnace blower. However, drainage holes in a blower housing do not prevent water from first entering the blower. The presence of water within the blower may cause deterioration or damage to blower components before the water eventually flows through the blower drainage holes.\nAccordingly, the invention described and claimed herein relates to a sleeve which provides a junction between a furnace blower pipe and a heating system pipe and which provides means for water condensation drainage away from a furnace blower. The sleeve comprises two generally concentric cylinders, an outer cylinder and an inner cylinder, which are joined by an annular wall. A blower pipe cavity which is formed between the two cylinders on one side of the annular wall is sized to mate with an end of a blower pipe. An exhaust pipe cavity which is formed between the two cylinders on an opposite side of the annular wall is sized to fit an end of a heating system exhaust pipe. The exhaust pipe cavity also serves as a condensate reservoir. At least one protrusion from the annular wall or the outside cylinder into the exhaust pipe cavity forms a pipe-stop which prevents the exhaust pipe from abutting the annular wall thereby providing a space between the annular wall and the exhaust pipe for water condensate to flow. At least one hole through the outer cylinder and into the exhaust pipe cavity provides a condensate exit port from the exhaust pipe cavity. It is a primary advantage of the invention to provide means to drain water from the inner exhaust pipe walls away from the blower and thereby prevent or significantly reduce the flow of water into a furnace blower.\nIt is a further primary advantage of the invention to provide a drainage outlet for water condensation that does not interfere with the operation of a furnace blower. It is still another primary advantage of the invention to provide a condensate evacuation apparatus which is adapted for use on horizontally, vertically or diagonally mounted pipe systems.\nThose skilled in the art will recognize that the invention is not limited to furnace pipe applications. The apparatus of the invention may be used to evacuate water condensation from a pipe at any number of similar pipe junctions. It is an advantage of the invention to provide a condensation drainage apparatus that can accommodate a wide variety of pipe sizes. These and other objects of the invention will be made apparent in the drawings and in the following detailed description of the invention."}
{"text": "1. Technical Field\nAspects of this document relate generally to air fresheners.\n2. Background Art\nAir fresheners are known to come in a variety of forms, including chemical compositions that are activated by the air, fan operated devices that push or pull a scent into a room, and cardboard type materials that provide a scent in the air. The cardboard type materials are similar to an automotive tree air freshener that emits a scent and the freshener's effectiveness steadily decreases over time until the air freshener is nothing more than a decorative element. Chemical compositions may include gels or similar compositions that are absorbed into the air to mask any odors. Fan operated devices may be battery or outlet powered and typically push or pull air across a scented portion of the device and into the surrounding area. The scented portion may be connected to a scented liquid for absorption and the fan speed may be adjusted to vary the strength of the scent within the room."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a method for forming a tungsten silicide layer, and more particularly, the present invention relates to a method for depositing a tungsten silicide layer (WSix) using a tungsten source gas and a silicon source gas containing chlorine.\n2. Description of the Related Art\nAs semiconductor devices become highly integrated, the size of patterns formed on individual chips become smaller and an interval between the patterns becomes narrower. These ultra-fine patterns create problems with the conventional use of polysilicon as a wiring material for a gate electrode and a bit line. That is, as the size of the patterns become smaller, the specific resistivity of polysilicon is increased, which in turn can result in an RC time delay and an IR voltage drop. For this reason, polycide, which has characteristics similar to those of polysilicon but with much smaller specific resistivity (e.g., several times to several tens times smaller) has been suggested as a wiring material for the gate electrode and the bit line of a VLSI circuit. More specifically, a composite layer of polysilicon and a refractory metal silicide has been used as a polycide of wiring electrodes.\nSilicides of refractory metals, such as tungsten (W), molybdenum (Mo), titanium (Ti) and tantalum (Ta), have been adapted for use as a low resistive wiring material suitable for manufacturing the VLSI circuit. Silicides are bonded to polysilicon doped with high-density impurities, whereby a gate electrode having a polycide structure is formed. Preferably, a low pressure chemical vapor deposition (LPCVD) method is used for depositing refractory metal silicides. Tungsten silicide, when bonded to polysilicon, exhibits superior characteristics with respect to self-passivation, stability against wet chemicals, surface roughness, adhesion, anti-oxidation and reproducibility.\nTungsten silicide (WSix) thin films are formed using monosilane (SiH4) and tungsten hexafluoride (WF6) as a precursor gas and are deposited on semiconductor substrates by LPCVD. However, certain drawbacks are associated with the LPCVD process. One such drawback is that tungsten silicide cannot be conformally deposited on a stepped portion of a substrate. Another drawback is that a substantial amount of fluorine remains in the deposited tungsten silicide thin film, which in turn can result in operational defects in the manufactured semiconductor device. That is, when the semiconductor wafer is later exposed to a temperature exceeding 850xc2x0 C. during an anneal process, the remaining fluorine ions move into a lower silicon oxide layer through a polysilicon layer. An effective thickness of the silicon oxide layer is thus increased, resulting in inconsistent electric characteristics of the semiconductor device having the silicon oxide layer.\nFor this reason, it has been suggested to deposit the tungsten silicide layer using dichlorosilane (DCS; SiH2Cl2) instead of monosilane. A tungsten silicide layer which has been deposited using dichlorosilane (hereinafter, referred to as a xe2x80x9cDCS tungsten silicide layerxe2x80x9d) exhibits superior step coverage and low fluorine content as compared with a tungsten silicide layer which has been deposited using monosilane (hereinafter, referred to as a xe2x80x9cMS tungsten silicide layerxe2x80x9d).\nIn view of its advantages, the DCS tungsten silicide layer has been preferred over the MS tungsten silicide layer as the wiring material for a word line (that is, gate electrode) and a bit line. However, the DCS tungsten silicide layer is not without its drawbacks. That is, various problems result from the fact that when the dichlorosilane gas is reacted with tungsten hexafluoride (WF6) gas, chlorine radicals remain on a surface and an inner portion of a thin film.\nFIG. 1 is a flow chart showing a conventional method for depositing a DCS tungsten silicide layer.\nReferring to FIG. 1, a silicon wafer, that is, a semiconductor substrate having a polysilicon layer at an uppermost region thereof is loaded in a process chamber of a plasma-enhanced chemical vapor deposition (PECVD) apparatus (step S1). Then, tungsten hexafluoride (WF6) gas and dichlorosilane (SiH2Cl2) gas are introduced into the process chamber to form a tungsten silicide nucleus on a surface of the polysilicon layer (step S3). Then, under reaction conditions of a pressure of about 1 to 1.5 Storr, a deposition time of 10 to 25 seconds, and a temperature of 450 to 700xc2x0 C., tungsten hexafluoride (WF6) gas and dichlorosilane (SiH2Cl2) gas are introduced into the process chamber in a ratio of about 13:180. As a result, a tungsten silicide layer is deposited on the surface of the polysilicon layer having the tungsten silicide nucleus (step S5). Then, a post-flushing process is carried out for about 10 seconds to reduce stress by introducing a monosilane (SiH4) gas into the process chamber (step S7) at flow rate of about 300 sccm.\nFIG. 2A is a graph showing a chlorine profile in the tungsten silicide layer and the polysilicon layer when the DCS tungsten silicide layer has been deposited, and FIG. 2B is a graph showing a chlorine profile in the tungsten silicide layer and the polysilicon layer after an annealing process has been carried out.\nGenerally, the DCS tungsten silicide layer is deposited at a temperature of 620xc2x0 C., which is about 200xc2x0 C. higher than a deposition temperature of the MS tungsten silicide layer. Also, the DCS tungsten silicide layer exhibits a columnar structure in which a hexagonal phase and a tetrahedral site co-exist. In addition, the DCS tungsten silicide layer has a body-centered cubic structure in which a covalent bond is formed between one tungsten atom and eight silicon atoms. Excess silicon, which is not involved in the covalent bond, acts as a stacking fault. Silicon, which acts as the stacking fault, is bonded to chlorine dissociated from dichlorosilane gas, thereby forming a silicon chloride crystal. Chlorine breaks Sixe2x80x94H and Sixe2x80x94Si bonds so as to be bonded to silicon and excess silicon is re-bonded to chlorine because of the bond strength differences. Therefore, as shown in FIG. 2A, chlorine is diffused from the tungsten silicide layer into the underlying polysilicon layer as the result of a domino phenomenon of continuous dissociation and crystal-bonding, so that chlorine in the form of SiClx crystals is concentrated on an interfacial surface between the tungsten silicide layer and the polysilicon layer.\nIn addition, after the anneal process has been carried out at a temperature of about 800xc2x0 C., the tungsten silicide layer has a stable stoichiometry due to the tetrahedral site crystalline structure, whereby an amount of chlorine considered as SiClx crystals increases in a region of the polysilicon layer which has a relatively sufficient amount of silicon (referred to FIG. 2B).\nAs described above, according to the conventional DCS tungsten deposition method, a substantial amount of SiClx crystal remains in the thin film and a surface of the polysilicon layer, whereby a Milky Way phenomenon (that is, a collection of twinkling points viewed in an optical microscope) occurs which is caused by a scattering of light due to the SiClx crystal, thereby resulting a visual defect. The visual defect affects subsequent photolithography processes, making it difficult to carry out such photolithography processes.\nIn addition, the chlorine atoms contained in the DCS tungsten silicide layer and the polysilicon layer diffuse into an inside of the silicide layer in the form of SiClx crystals, and as a result, the polysilicon layer is abnormally grown.\nThat is, while the crystallization temperature of the polysilicon layer is about 530 to 550xc2x0 C., the polysilicon layer is initially deposited in an amorphous state and crystallization of the polysilicon layer is carried out in the process chamber at a temperature of about 620xc2x0 C. when the DCS tungsten silicide layer is deposited. At this time, excess silicon contained in the DCS tungsten silicide layer is bonded to chlorine to form the SiClx crystal. In this state, the SiClx crystal diffuses from the tungsten silicide layer into the underlying polysilicon layer as a result of the domino phenomenon of continuous dissociation and crystal-bonding. Due to the SiClx crystals, silicon dissociated from Sixe2x80x94Si bond or Sixe2x80x94H bond diffuses towards the polysilicon layer, thereby causing growth and crystallization of the polysilicon layer. In particular, the monosilane post-flushing process, which is carried out after the DCS tungsten silicide layer is deposited, continuously supplies excess silicon into the tungsten silicide layer. Therefore, after the DCS tungsten silicide layer deposition process and the monosilane post-flushing process have been finished, the growth and crystallization of the polysilicon layer in the amorphous state occurs locally from the diffusion of excess silicon. Accordingly, the polysilicon layer protrudes through the tungsten silicide layer formed thereon, thereby resulting a crack known as a xe2x80x9chaze phenomenonxe2x80x9d.\nAccording to one aspect of the present invention, there is provided a method of forming a tungsten silicide layer, which includes loading a semiconductor substrate having a polysilicon layer into a process chamber of a plasma enhanced chemical vapor deposition device, and introducing a silicon source gas, a tungsten source gas, and a hydrogen compound gas for reducing a chlorine radical into the process chamber, to thereby deposit the tungsten silicide layer on the polysilicon layer.\nAccording to another aspect of the present invention, there is provided a method for forming a tungsten silicide layer, which includes loading a semiconductor substrate having a polysilicon layer into a process chamber of a plasma enhanced chemical vapor deposition device, pre-flushing a surface of the polysilicon layer by introducing a hydrogen compound gas for reducing a chlorine radical into the process chamber, introducing a silicon source gas, a tungsten source gas and the hydrogen compound gas into the process chamber so as to form a tungsten silicide nucleus on the surface of the polysilicon layer, and introducing the silicon source gas and the tungsten source gas into the process chamber so as to deposit the tungsten silicide layer on the polysilicon layer.\nAccording to still another aspect of the present invention, there is provided a method for forming a tungsten silicide layer, which includes loading a semiconductor substrate having a polysilicon layer into a process chamber of a plasma enhanced chemical vapor deposition device, pre-flushing a surface of the polysilicon layer by introducing a hydrogen compound gas for reducing a chlorine radical into the process chamber, introducing a silicon source gas, a tungsten source gas and the hydrogen compound gas into the process chamber so as to form a tungsten silicide nucleus on the surface of the polysilicon layer, and introducing the silicon source gas, the tungsten source gas, and the hydrogen compound gas into the process chamber so as to deposit the tungsten silicide layer on the polysilicon layer.\nAccording to the present invention, when the tungsten silicide layer is deposited on the polysilicon layer by introducing, for example, a tungsten hexafluoride (WF6) gas and a dichlorosilane (Si2HCl2) gas into the process chamber, a hydrogen compound gas such as B2H6, PH3, AsH3 and/or NH3 is also introduced into the process chamber. As a result, a chlorine radical dissociated from the dichlorosilane gas by the hydrogen compound is reduced into hydrogen chloride (HCl) and is removed together with an exhaust gas. Accordingly, visual defects caused by SiClx crystals and abnormal growth of the polysilicon layer may be avoided."}
{"text": "Magnetic resonance imaging (MRI) is currently contra-indicated for patients who have implanted medical leads. This is due largely to the patient safety issue that results when the strong electromagnetic fields of an MRI system interact with the antenna-like therapy delivery leads of an active implantable medical device (AIMD). It is well documented that the radio frequency (RF) signals that are generated by the MRI system can couple along the length of a lead conductor or conductors and create induced RF currents. These RF currents can cause significant heating at points of high current concentration, the most significant of which is at the distal tip electrode, where the lead system makes direct contact with body tissue. Excessive RF currents at the point of electrode contact to tissue can create a serious or even life-threatening situation. In the literature, there are reports of MRI induced thermal damage to brain tissue from around the area of implanted deep brain stimulation electrodes, pacemaker pacing capture threshold changes or even loss of capture which means that a pacemaker dependent patient is without life support.\nAIMDs can include completely implantable systems or a combination of externally worn devices with implanted leads. AIMDs include the group of cochlear implants, piezoelectric sound bridge transducers and the like. AIMDs can also include a variety of neurostimulators and brain stimulators. Sometimes these are called neuromodulators. Neurostimulators are used to stimulate the vagus nerve, for example, to relieve epilepsy, obesity, and depression. Deep brain stimulators are pacemaker—like devices and include electrodes implanted deep into the brain matter for sensing the onset of the seizure and also providing electrical stimulation to brain tissue to prevent the seizure from actually occurring. There are also other types of deep brain stimulators used to correct other abnormalities, such as Tourette's Syndrome and the like. AIMDs also include all types of cardiac pacemakers, left ventricular assist devices (LVADs), artificial hearts and implantable cardioverter defibrillators. AIMDs also include drug pumps which can be used for dispensing of insulin, chemotherapy drugs, pain medications and the like. AIMDs can include a variety of implanted or external-implanted bone growth stimulators for rapid healing of fractures. AIMDs may also include urinary incontinence devices, pain release spinal cord stimulators, anti-tremor stimulators, congestive heart failure devices, cardiac resynchronization therapy devices (CRT), and the like.\nAs disclosed in US 2007/0112398 A1 and U.S. Pat. No. 7,853,325, the contents of which are incorporated herein by reference, a novel method to n minimize the expected heating at the distal tip of the lead system is to incorporate a bandstop filter. This bandstop filter is comprised of an inductor and capacitor in parallel, with the bandstop filter connected in series with one or more conductors of the implanted lead system. In such a system, the bandstop filter is constructed so that its resonant frequency or frequencies coincides with the RF operating frequency of one or more MRI systems.\nRF frequencies are directly related to the MRI machine static magnetic field by the Lamour Relationship wherein the frequency is equal to 42.56 times the static field strength in Teslas (for hydrogen scanners). Typical MRI RF pulsed frequencies are 64 MHz for 1.5 T systems and 128 MHz for 3.0 T systems. At resonance, the impedance of the bandstop filter is quite high (for example, above 2,000 ohms) which reduces the flow of distal electrode to tissue current at the MRI RF pulsed frequency thereby reducing implanted lead and/or electrode heating. Increasing the impedance at the lead distal tip electrode greatly reduces the amount of RF current that would flow into body tissue. It has been documented that excess current can cause tissue damage, pacing capture threshold (PCT) changes or even tissue necrosis.\nImplementation of this technology in implantable leads is a significant challenge. Bandstop filters for use in implantable lead systems must be biocompatible, not significantly change the electrical performance characteristics of the lead (except within the context of the invention), and must not significantly affect size, weight, or implantability. With increasingly smaller leads being developed to accommodate small vasculature and left ventricular pacing through the coronary sinus, bandstop technology must be equally scalable to match the same demands.\nThe present invention satisfies these needs and provides other related advantages."}
{"text": "FIG. 1 is a drawing illustrating a semiconductor device in the form of a “system in a package” (SiP) obtained through a related art semiconductor manufacturing method.\nReferring to FIG. 1, a related art SiP semiconductor device may include interposer 11, first device 13, second device 15, and third device 17.\nFirst device 13, second device 15, or third device 17 may include any one of a central processing unit (CPU), a static random access memory (SRAM), a dynamic random access memory (DRAM), a flash memory, a logic large scale integrated circuit (LSI), a power integrated circuit (IC), a control integrated circuit (IC), an analog large scale integrated circuit (LSI), a microwave monolithic integrated circuit (MM IC), a complementary metal-oxide semiconductor radio-frequency integrated circuit (CMOS RF-IC), a sensor chip, a micro-electro-mechanical systems (MEMS) chip, etc.\nConnection units may be formed between first and second devices 13 and 15, and between second and third devices 15 and 17, and may transfer signals to each device.\nHowever, to extensively use the SiP semiconductor device having the above structure, a problem regarding heat dissipation may need to be solved. Particularly, the problem regarding heat dissipation of a device, such as second device 15 positioned in a mid-layer, may need to be solved to extensively use the SiP semiconductor device."}
{"text": "1. Field of the Invention\nThis invention relates generally to improvements in toilet paper dispensing apparatus.\n2. Description of the Prior Art\nVarious constructions are known for supporting a roll of toilet paper in a conventional fashion, that is wherein the axis of the roll is substantially parallel to a wall or other supporting structure within a bathroom. Representative of one such construction are the patents to Pendergast, U.S. Pat. No. 2,579,201 and Pinkham, Jr., U.S. Pat. No. 2,873,158. In each of these instances sidewalls are provided for supporting the ends of the roll and the apparatus includes a hinged top or front portion which is movable between open and closed positions.\nOther representative patents are those to Patterson et al., U.S. Pat. No. 2,487,763 and Marchand, U.S. Pat. No. 2,576,526. In each of these patents, the invention incorporates a casing or enclosure having a pivotally mounted front portion supporting a roll and movable between open and closed position.\nIn the patent to Chevas, U.S. Pat. No. 3,750,971, a cantilevered roll extends outwardly from one sidewall of the casing or enclosure and the apparatus further includes an integral top and opposite sidewall hinged for movement between open and closed positions.\nStill another representative of the prior art is the patent to Girard U.S. Pat. No. 3,475,067, which discloses a dispenser intended for mounting to the upper rim of the toilet bowl rather than to a wall or partition of the bathroom."}
{"text": "1. Technical Field\nThis invention relates to diode-pumped lasers and more specifically to configurations of resonator cavities for solid-state lasers.\n2. Discussion of the Prior Art\nLaser diode light emissions are commonly dedicated to \"pumping\" gain materials to achieve lasing in larger host cavities which in turn output continuous waves or pulses of light. Laser diode side-pumping complicates controlling transverse modes of lasing in solid-state gain materials such as 1% neodymium-doped yttrium aluminum garnet (Nd:YAG) or yttrium lithium fluoride (Nd:YLF), but end-pumping facilitates controlling transverse modes. Laser diode light pumped into a solid-state gain material is converted along the light's \"absorption length\" (typically four to five mm) most efficiently within an axially cylindrical TEM.sub.00 mode volume through the gain material. Solid-state laser gain \"G\" is inversely proportional to the sum of the cross-sectional areas of the pumping beam A.sub.p and of the cavity mode A.sub.c : ##EQU1##\nHigher power laser diodes have light-emitting junctions or \"slits\" which are dimensioned non-symmetrically, typically one micron high vertically and many times as long horizontally. Laser diodes pump light in beams which diverge non-symmetrically, typically ten degrees horizontally and fifty degrees vertically full-width at half-maximum intensity (\"FWHM\"). Any laser diode pump light beams which diverge more widely than ten degrees beyond their \"waists\" spread, through several millimeters, to objectionably broad spots. Fortunately, a laser diode beam initially only one micron high can afford to have its short vertical height increased by magnification five times in a trade-off to achieve having its excessive vertical divergence decreased by collimation from fifty degrees inward to ten degrees, to be suitable for imaging axially into a TEM.sub.00 mode volume.\nNon-symmetrically divergent elliptical beams having their collimations improved by vertical magnification fortuitously tend to become rounder. Minimally elliptical beams (i.e., those with short, say.ltoreq.several hundred micron, lengths not dominating their heights) undergoing vertical magnification become fairly circular. Diode pump beam-rounding distributes light energy more uniformly, and thereby reduces thermal focusing and gain-guiding effects in pumped mode volumes of solid-state laser cavities, which in consequence reduces distortions in output beams.\nLaser output beams are preferably circular for applications in general. The transverse shape of a laser output beam corresponds to that of the lasing mode beam inside the cavity at the out-coupling mirror. This predisposes solid-state laser designers to specify cavities for resonating in cross-sectionally circular modes. To match circular mode areas, pumping beams are preferably also circular in cross-section.\nNon-symmetrically divergent light beams may be re-shaped for example by passage through a cylindrical lens or a prism, or by reflections against a pair of concave mirrors oriented off-axis as described in U.S. patent application Ser. No. 185,466, filed Apr. 25, 1988 and assigned to the assignee of this application. Optical cavities have also used prisms for different purposes, for example tilted steeply in \"An Improved Line Narrowing Technique for a Dye Laser Excited by a Nitrogen Laser,\" Optical Communications vol. #4 pg. 187 (1971).\nLaser output beams have been made more powerful conventionally by cavities being pumped with laser diodes in turn made more powerful through their junctions being lengthened proportionately. Laser diode powers and operational lifetimes also have been further increased by arranging multiple arrays side-by-side to form long \"bars\". For example, thirty discrete one hundred micron-long diode arrays are aligned on 333 micron centers in ten-watt bar model SDL 3490-S from the Spectra Diode Labs. Such a diode-array bar has a slit with a very high aspect ratio of one micron high by one centimeter long.\nUnfortunately however, laser-diode radiations are not diffraction-limited, i.e. longer diode pump beams do not have narrower divergences. An optical system which images a laser diode pump beam by horizontally reducing its length simultaneously increases its divergence so that before the pumping light can be entirely absorbed, the beam expands beyond the diameter of the fundamental (TEM.sub.00) mode area in a solid-state gain medium. This results either in lasing in higher transverse modes or, in cavities restricted to the fundamental mode (for example by an aperture), less efficient lasing.\nDespite thermal effects and gain-guiding considerations, thin elliptical laser diode beams can be used to pump solid-state lasing gain cavities. U.S. Pat. No. 4,653,056 by Baer describes multi-strip arrays of diodes which emit light compiled into beams possessing too much spacial structure and having poor focusing qualities (col. 2 lines 52-61). This disadvantage may be addressed by configuring a cavity to expand its lasing mode diameter to match the long diameter of the focused image of the pump beam, which allows utilizing longer diode bars with higher powers. However, Baer considers only circular modes.\nIncreases in diode pump power accompanied by proportional increases in diode length dictate equal increases in diameters of circular modes in end-pumped solid-state lasers. Since solid-state laser gain G is inversely proportional to the cross-sectional area A.sub.c of the cavity mode volume, diode pump power-increases accompanied by length-increases are traded off against gain decreases. This poses dilemmas in making longer diode arrays or array-bars to be compatible for end-pumping conventional solid-state laser cavities.\nLaser diodes in discrete arrays arranged either serially along a bar or individually can be used to pump light into ends of optical fibers, which can in turn either separately couple the light into multiple-bounce point cavities (see Optics Letters, Vol. 13, No. 4, pg. 306, April 1988), or be bundled together to couple the light into single- or multiple-end cavities.\nU.S. Pat. No. 4,785,459 by Baer describes multiple laser diode arrays each two hundred microns long centered one mm apart along a bar. The bar must be precisely and meticulously aligned to couple the light from the diodes into respective bounce points also one m apart along a multi-bounce solid-state gain element of Nd:YAG.\nThere remains, therefore, a need for a laser cavity configuration to facilitate converting longer laser diode pump light beams efficiently and conveniently into TEM.sub.00 lasing output energy."}
{"text": "This invention relates to a new and useful improvement relating to the treatment of athletic type socks so as to prevent irritation and blisters of the feet.\nIt is a well known fact that active persons frequently have foot irritations from shoes when engaged in athletics or other activities, particularly in initial or training portions of sport programs. This irritation of the feet is primarily caused by the relative movement of the foot within the sock in that such relative movement causes friction which usually results in discomfort and irritation and even incapacitating blisters. Many medications have been provided to treat irritated feet but such medications are only intended to give relief to the damage already done to feet rather than providing a preventative to the actual source of the trouble."}
{"text": "1. Field of Invention\nThe invention relates to instruments for drawing an ellipse and more particularly to an improved instrument for drawing an ellipse or a closed line having three or more arcuate portions.\n2. Description of Related Art\nThere have been numerous suggestions in prior patents for instruments for drawing an ellipse. For example, U.S. Pat. No. 4,182,043 (“the '043 patent”) describes a drawing instrument for an ellipse. However, the '043 patent is limited in drawing different shapes of ellipse. Thus, continuing improvements in the exploitation of instrument for drawing ellipse or the like are constantly being sought."}
{"text": "The present invention relates to a correction method of document inclination angle for an image reading apparatus such as a scanner, facsimile device, digital copier in which a document is placed on an image reading position, and an image of the document is read, and specifically to an image reading apparatus in which, even when the placed document is inclined, it can be corrected.\nIn an image reading apparatus such as a scanner, facsimile device, digital copier in which a document is placed on an image reading position, and an image of the document is read, there is a case in which the document is placed while being inclined. When the inclined document is read, for example, copied (image formation), as it is, there is a possibility that some portions of the image are lost.\nAccordingly, in the conventional image reading apparatus, an inclination angle of the document is detected by detecting an inclination of a row of characters from image information obtained from the inclined document, and the inclination angle is corrected corresponding to obtained image information.\nIn the above method, in order to detect the inclination angle of the document, it is necessary that image information is read (pre-scanned) after the document is previously set on a reading position. Therefore, a long period of time is necessary for detecting the inclination angle, and thereby, image reading efficiency is lowered. Further, in order to detect the inclination of a character row, it takes a long time because complicated calculation is necessary, and further, cost is increased because many memories are used for calculation."}
{"text": "One of the challenges facing every community is the collection and disposal of refuse, such as household and commercial garbage. Generally, the refuse is ultimately disposed of at a permanent, sanitary landfill location, although certain types of garbage may be incinerated. It is desirable to provide for methods and systems that tidily, efficiently and economically collect and transport the refuse to the ultimate disposal site.\nIn densely populated urban areas, relatively small collection vehicles have traditionally collected refuse from individual house-holds and dwelling-places and directly transported the collected refuse to a nearby landfill site. This is economically inefficient because of the relatively small capacity of the collection vehicles, resulting in a relatively high cost of transporting a given weight of refuse.\nRemote rural communities face added challenges because they are typically nowhere near a landfill site. Accordingly, it is often necessary for these communities to transport their collected refuse to far-away landfill sites. In order to maximize the cost efficiency of this added burden, the refuse transportation is often done via large transportation trucks which have a much greater capacity than the aforementioned collection vehicles. It is thus important to provide a location in or around remote communities where refuse can be collected, temporarily stored and ultimately transferred to the large trucks for its ultimate haul to a landfill site.\nIt has proven difficult to address these challenges in an effective manner. Canadian Patent No. 1,132,329, issued Sep. 28, 1982 to Neufeldt, represents one attempt. Neufeldt discloses a refuse disposal method and apparatus which utilizes a transfer container adapted to receive refuse from collection vehicles and discharge the refuse into larger transportation vehicles. Neufeldt accomplishes this by providing the transfer container mounted alongside a steep embankment so that the collection vehicles can approach same via a gentle ramp and backfill the hopper portion of the container. Ultimately, Neufeldt's container is pivoted, resulting in a \"tilted\" dumping of the refuse into a transportation vehicle located below the container and adjacent the steep embankment.\nNeufeldt, however, has inherent disadvantages. First, the physical landscape required to accommodate the container assembly is impractical. Naturally-occurring locations offering a steep embankment of such a height and configuration to accommodate the Neufeldt assembly are rare. Accordingly, it will usually be necessary to construct a steep embankment through piling and shaping dirt, gravel or the like to a suitable height and providing suitable roads to accommodate the collection and transportation vehicles. Not only is this process extremely expensive, but it is also environmentally-unfriendly since it will be very difficult to reclaim the site should the collection container be moved in the future.\nSecond, Neufeldt does not offer the users the ability to regulate the amount of refuse \"dumped\" from the container to the transportation truck. Instead, the container is simply \"tilted\", thereby dumping the entirety of its contents. This is acceptable for transportation trucks which have the capacity for all of the collected refuse. However, transportation trucks may arrive at a particular remote community partially full because the same truck services numerous communities in the same trip. Thus, if a transportation truck already contains sufficient refuse so that it cannot accept the entirety of the contents of the Neufeldt container, the user is faced with one of two undesirable choices. First, he can bypass accepting the refuse at that time and simply collect it on another trip. This is disadvantageous because: (a) the transportation vehicle would be returning with less than 100% capacity, thus descreasing the cost-effectiveness of the operation; and (b) the remote community may not have the capacity to store the collected refuse while awaiting a second trip sometime later. The alternative to refusing to accept the refuse is to simply dump the container and let the excess spill over the top of the transportation vehicle. This is obviously disadvantageous since the spillage would require clean-up, or if it were not cleaned up, would be unsightly and would produce unpleasant odours and attract scavenging animals.\nAccordingly, the present invention provides an efficient, cost-effective and self-contained apparatus and a method that allows a user to regulate the amount of collected refuse transferred to a transportation vehicle in any given dump."}
{"text": "Conventional paper shredding machines mostly include a roller blade set constructed of two roller blades that shred or cut paper to be fed into strips or scraps as a result of the opposed rotations of the two roller blades such that information as recorded on the paper is destroyed for confidentiality, and the strips or scraps of paper can be easily compressed to reduce processing space. However, optical discs, regardless of CD-ROM discs or CD-R discs, rather than paper have evolved to be one of the major means for storing information. Once information contained in such optical discs has lost its original value and needs to be destroyed, manually breaking the optical discs not only cannot destroy the information as stored, it also may cause personal injuries. An optimum measure is to shred these discs by means of mechanical operations such as those in conventional paper shredding machines.\nIn the highly economized society as of today, plastic money, such as credit cards, debit cards, ATM cards, or even membership cards issued by enterprises for promotional purposes, and registration cards issued by medical institutions, have made “cards” become an article that can certainly be found in everyone's pocket. When these cards have expired or been replaced with new cards, the most commonly adopted measure is to cut the cards in halves for disposal. However, danger still exists in such a disposing measure because most cards carry the user's signature and the registration cards may also carry personal, medical history, or personal information. It is possible that other individuals with malicious intention may still have access to these halved cards."}
{"text": "The virtual reality (VR) technology is a computer simulation system, which can generate a simulated environment and provide an immersive experience for a user by way of three-dimensional dynamic scenes, interactions among entities and so on.\nIn related conventional technologies, a simple and inexpensive virtual reality solution is provided by providing virtual reality glasses including virtual reality function components and by using an external electronic device such as a user's mobile phone or tablet to play virtual reality content.\nHowever, different external electronic devices may have big differences in sizes, specifications, materials and so on. If an external electronic device is not effectively restricted at a position, the external electronic device is prone to vibrate, shake, or the like during use, which affects a user's immersive experience during use, and even causes the user to feel discomfort such as dizziness.\nThe virtual reality glasses of the present disclosure are directed towards overcoming one or more the problems set forth above."}
{"text": "Round balers for producing round bales may include a number of rotatable roller assemblies about which bale-forming belt(s) may be trained. An arrangement for mounting such a roller assembly is disclosed in U.S. Pat. No. 6,623,167, the entirety of which is hereby incorporated by reference herein."}
{"text": "1. Field of the Invention\nThe present invention concerns advertising. In particular, the present invention concerns improving advertising using end user local time information, such as local time-of-day, local day-of-week, local date, and/or local season information (referred to generally as “time information” in the specification without loss of generality).\n2. Related Art\nAdvertising using traditional media, such as television, radio, newspapers and magazines, is well known. Unfortunately, even when armed with demographic and entirely reasonable assumptions about the typical audience of various media outlets, advertisers recognize that much of their ad budget is simply wasted. Moreover, it is very difficult to identify and eliminate such waste.\nRecently, advertising over more interactive media has become popular. For example, as the number of people using the Internet has exploded, advertisers have come to appreciate media and services offered over the Internet as a potentially powerful way to advertise.\nAdvertisers have developed several strategies in an attempt to maximize the value of such advertising. In one strategy, advertisers use popular presences or means for providing interactive media or services (referred to as “Websites” in the specification without loss of generality) as conduits to reach a large audience. Using this first approach, an advertiser may place ads on the home page of the New York Times Website, or the USA Today Website, for example. In another strategy, an advertiser may attempt to target its ads to narrower niche audiences, thereby increasing the likelihood of a positive response by the audience. For example, an agency promoting tourism in the Costa Rican rainforest might place ads on the ecotourism-travel subdirectory of the Yahoo Website. An advertiser will normally determine such targeting manually.\nRegardless of the strategy, Website-based ads (also referred to as “Web ads”) are often presented to their advertising audience in the form of “banner ads”—i.e., a rectangular box that includes graphic components. When a member of the advertising audience (referred to as a “viewer” or “user” in the Specification without loss of generality) selects one of these banner ads by clicking on it, embedded hypertext links typically direct the viewer to the advertiser's Website. This process, wherein the viewer selects an ad, is commonly referred to as a “click-through” (“Click-through” is intended to cover any user selection). The ratio of the number of click-throughs to the number of impressions of the ad (i.e., the number of times an ad is displayed or otherwise rendered) is commonly referred to as the “click-through rate” or “CTR” of the ad.\nA “conversion” is said to occur when a user consummates a transaction related to a previously served ad. What constitutes a conversion may vary from case to case and can be determined in a variety of ways. For example, it may be the case that a conversion occurs when a user clicks on an ad, is referred to the advertiser's Web page, and consummates a purchase there before leaving that Web page. Alternatively, a conversion may be defined as a user being shown an ad, and making a purchase on the advertiser's Web page within a predetermined time (e.g., seven days). In yet another alternative, a conversion may be defined by an advertiser to be any measurable/observable user action such as, for example, downloading a white paper, navigating to at least a given depth of a Website, viewing at least a certain number of Web pages, spending at least a predetermined amount of time on a Website or Web page, etc. Often, if user actions don't indicate a consummated purchase, they may indicate a sales lead, although user actions constituting a conversion are not limited to this. Indeed, many other definitions of what constitutes a conversion are possible. The ratio of the number of conversions to the number of impressions of the ad (i.e., the number of times an ad is displayed or otherwise rendered) is commonly referred to as the conversion rate. If a conversion is defined to be able to occur within a predetermined time since the serving of an ad, one possible definition of the conversion rate might only consider ads that have been served more than the predetermined time in the past.\nThe hosts of Websites on which the ads are presented (referred to as “Website hosts” or “ad consumers”) have the challenge of maximizing ad revenue without impairing their users' experience. Some Website hosts have chosen to place advertising revenues over the interests of users. One such Website is “Overture.com,” which hosts a so-called “search engine” service returning advertisements masquerading as “search results” in response to user queries. The Overture.com Website permits advertisers to pay to position an ad for their Website (or a target Website) higher up on the list of purported search results. If such schemes where the advertiser only pays if a user clicks on the ad (i.e., cost-per-click) are implemented, the advertiser lacks incentive to target their ads effectively, since a poorly targeted ad will not be clicked and therefore will not require payment. Consequently, high cost-per-click ads show up near or at the top, but do not necessarily translate into real revenue for the ad publisher because viewers don't click on them. Furthermore, ads that viewers would click on are further down the list, or not on the list at all, and so relevancy of ads is compromised.\nSearch engines, such as Google for example, have enabled advertisers to target their ads so that they will be rendered in conjunction with a search results page responsive to a query that is relevant, presumably, to the ad. Although search result pages afford advertisers a great opportunity to target their ads to a more receptive audience, search result pages are merely a fraction of page views of the World Wide Web.\nSome online advertising systems may use ad relevance information and document content relevance information (e.g., concepts or topics, feature vectors, etc.) to “match” ads to (and/or to score ads with respect to) a document including content, such as a Web page for example.\nThe foregoing ad serving systems can be thought of as keyword-targeted systems (where ads are targeted using terms found in a search query) and content-targeted systems (where ads are targeted using content of a document). Although keyword-targeted and content-targeted ad systems have improved the usefulness of ads, and consequently their performance (e.g., in terms of click-through rate, conversion rate, etc.), there is still plenty of room for improvement. Such improvement can be expected with better targeting.\nSome ad servers (e.g., bid management tools) allow advertisers to specify time of day ad targeting. This offers some improvement in ad targeting. Unfortunately, however, time-of-day ad targeting techniques use a fixed time zone in their ad serving determination. Even if such businesses have ads that are relevant to a search query or a Web page, if the end user viewing a search results Web page or the content of a Web page is outside a preferred temporal window, the ads may not be very useful and may not perform well.\nConsider, for example, a nationwide courier service that offers package pickup until midnight for next day delivery, while most of their competitors only offer package pickup until 8:00 PM, local time. If such a nationwide courier service wants to advertise on the Internet, the value to them of serving their ads from 7:30 PM to 11:55 PM (or some other time range corresponding to a time when users will not be able to meet an 8:00 PM pickup deadline, but should be able to meet a 12 midnight pickup deadline) might be higher than the value of serving their ads from 9:00 AM to 7:29 PM (because they have less competition after 8:00 PM and customers are likely to use their services that day). However, the value to the courier of serving its ads after 12:00 midnight might drop precipitously since the courier can no longer differentiate its services from its competitors for an end user that needs something to be delivered that day. Note in this example that the value of serving the ads depends on the time for the end user to which the ad will be rendered. For example, an ad touting pickup until midnight may have little value for a user on the east coast at 1:30 AM, but great value for a user on the west coast at 10:30 PM. If, however, the time-of-day is determined merely using a fixed time zone, this great difference in value cannot be used in ad serving.\nAs another example, advertisers may want to apply audience demographics which are relevant only if the local time of the audience is known. For example, an advertiser may want to increase the score of their ad during lunchtime on weekdays if demographic studies indicate that many people use the Internet to look for personal items during their lunchtime. Suppose an ad server (or a bid manager) is located in the Mountain time zone. At 10:30 AM Mountain time, it is only 9:30 AM Pacific time, but 12:30 PM Eastern time. Thus, at 10:30 AM Mountain time, the advertiser may want to increase the score of its ad for an end user on the east coast, but not for an end user on the west coast.\nIn view of the foregoing, there is a need for improving the usefulness, and consequently the performance, of advertisements. In particular, there is a need to allow businesses to better target their ads to a responsive audience."}
{"text": "1. Field of the Invention\nThe present invention pertains to the field of optoelectronics. The invention more particularly concerns an optoelectronic device which incorporates parallel optics and a three dimensional circuit board.\n2. Discussion of the Background\nOptoelectronic devices such as transceivers, transmitters, and receivers are known in the art. Furthermore, optoelectronic devices are known to incorporate parallel optics. Parallel optics are integrated into optoelectronic devices so as to increase the data density for a given unit of area or footprint and space. As an example, an eight channel optoelectronic device has eight photoactive devices. All eight of the photoactive devices can be transmitting devices such as a vertical cavity surface emitting laser (VCSEL) or a light emitting diode (LED) or some other light emitting device, or all eight of the photoactive devices can be receiving devices such as a photodetector or a PIN diode, or the eight photoactive devices can be a combination of light emitting devices and photodetectors.\nParallel optics optoelectronic devices contain electronics for modulating the photoactive elements. Typically the electronics are placed within or on the chassis of the optoelectronic device. However, as the modulation frequency is increased it is desirable to place the modulation electronics as close as possible to the photoactive devices. Given traditional board layout, a limit is reached as to how small the optoelectronic device can become due to the placement of the large number of electronics required while maintaining adequate heat dissipation.\nThus, there is a need for a device which is more compact (higher data density), more reliable, less costly, and easier to assemble than prior art devices."}
{"text": "As a storage apparatus, for example, a device (e.g., a disk array device) comprising a number of physical storage devices (e.g., HDDs (Hard Disk Drives)) is well known. This type of storage apparatus generally provides a logical volume to a host device. The host device, by transmitting an access requirement (a write requirement or a read requirement) specifying a logical volume, can write and read data to and from the logical volume.\nIn the Patent Literature 1, the computer described below is disclosed. That is, a logical volume (hereinafter referred to as an internal volume) is created in accordance with the disks in the computer, and the internal volume is provided to an application program. The data written from the application program to the internal volume is transferred to a storage apparatus outside the computer."}
{"text": "The present invention relates to tissue sampling or biopsy devices, and in particular to an automatic tissue sampling apparatus for utilizing two piece biopsy needle systems. In a specific aspect, the invention concerns a needle hub assembly for use with the automatic tissue sampling apparatus.\nA typical biopsy needle system may be utilized in obtaining tissue samples from a patient which includes two coaxial needles. The first needle generally consists of a substantially solid shaft in the form of a stylet having a cutting point at its end to facilitate insertion of the needle into the tissue to be sampled. The second needle is typically hollow in the form of a cannula and is disposed about the first inner needle.\nIn manual biopsy needle systems, the two needle components each have handles at their ends for manual operation. To obtain the tissue sample, the biopsy needle is inserted into the patient into the specific target tissue, preferably at the center of a desired cell mass to be investigated. The two needles are then reciprocated relative to each other to trap tissue within the outer needle cannula. Typically, the outer needle is held stable while projecting the first inner needle or stylet forward a short distance to penetrate the tissue. Then with the outer needle held stable, the inner stylet is retracted, severing the adjacent tissue and confining the tissue within the outer cannula. The entire biopsy needle assembly can then be withdrawn and the tissue sample sent for pathological examination.\nIn order to facilitate the tissue sampling operation, automatic tissue sampling or biopsy apparatus have been developed. One such device is depicted in U.S. Reissue Pat. No. 34,056 to Lindgren et al. This patent shows a tissue sampling device in which the two needles of the biopsy needle system are driven by a series of compressed springs. Each of the needles is supported by a hub, which hub is removably mounted within a respective slide or carriage. The actuating springs operate on the carriages to drive the stylet into the target tissue and then the outer cannula over the stylet to extract the tissue sample.\nAnother automatic tissue sampling device is shown in U.S. Pat. No. 5,284,156 to Shramm et al. This patent shows a biopsy apparatus that utilizes a different driving and actuation mechanism from that shown in the prior Lindgren patent. However, as with most automatic biopsy devices, the sampling apparatus in the Shramm patent also utilizes a carriage assembly to support hubs within which each of the needles are engaged.\nThe details of automatic tissue sampling or biopsy apparatus of this type can be easily gleaned from either the Lindgren or the Shramm patents. In accordance with the present invention, a tissue sampling apparatus 10, such as shown in FIGS. 1, 2, is provided which is similar in operation to the device shown in the Shramm U.S. Pat. No. 5,284,156. This biopsy device 10 includes a drive unit 11 which supports a biopsy needle assembly 12. As described above, the biopsy needle assembly 12 includes an outer cannula 13 and a concentrically disposed stylet 14. The stylet defines a tissue recess 15 which is used to capture and excise target tissue. The drive unit 11 includes a drive unit housing 17 that is provided with a hinged cover 18 to permit access to the internal workings of the drive unit 11. Many of the internal components of the drive unit are shown in FIG. 2. It is understood that the representation of tissue sampling device 10 in FIGS. 1, 2 is simply for purposes of illustration and that the device includes other components that are not depicted in the figures for clarity.\nThe particular tissue sampling device 10 shown in these figures can be a device sold as the MBD (Multiple Biopsy Device).TM. manufactured and marketed by Engineered Medical Systems, Inc., of Indianapolis, Ind. The details of the structure and operation of the MBD.TM. can be easily discerned from the product itself. However, for purposes of illustration in connection with the present invention, a general description of the components of the sampling device 10 will follow. In particular, the drive unit housing 17 is closed not only by the hinged cover 18 but also by a front end plate 19 and an opposite end plate 21. The front end plate 19 includes a slot 20 defined therein through which a biopsy needle assembly 12 extends. A support frame 22 is mounted within the drive unit housing 17 and includes a carriage stop 23 which restricts the longitudinal motion of the carriages within the unit.\nThe cannula and stylet of the biopsy needle assembly 12 are supported within the sampling device 10 by a cannula carriage 24 at the forward end of the device and a rear stylet carriage 25. The drive mechanism of the device includes a series of springs. A return spring 27 operates to provide a resistive force against reciprocation of the cannula carriage 24. A nest spring 28 and pusher spring 29 provide the actuation force to drive the two carriages 24, 25 forward into the tissue site. A center tube 30 extends through the cannula carriage 24 to maintain the coaxial alignment of the carriages and to support the shaft portion of the actuator handle 32. The actuator handle operates to cock the carriages 24, 25 into their operative positions. A latch 31 engages the carriages to hold them in their loaded or energized positions until the latch is released.\nAs thus illustrated, the tissue sampling device 10 is similar to automatic devices currently available in the marketplace. With these devices, the biopsy needle assembly 12 can be removed and sterilized and re-inserted into the same drive unit. Alternatively, a pre-sterilized disposable biopsy needle can be provided that fits into the carriages 24, 25 of the device. However, one difficulty with this approach is that the needle hubs affixed to the coaxial needles are often small and difficult to manipulate. In addition, the physician may have to spend time manipulating the hubs to properly position them within the carriages of whatever automatic biopsy device is being used. Finally, another difficulty encountered with these devices utilizing replaceable biopsy needle assemblies is that the assemblies are not easy to remove once a tissue sample has been obtained. Simple and safe removal is essential to protect the physician and to protect the excise tissue sample.\nThere is a need in the industry for a hub for use with a removable and replaceable biopsy needle assembly that is ergonomically efficient and is economical to produce. In addition, the hub should permit easy placement within the biopsy sampling device, and particularly in the carriages of the device."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention relate to the processing of metallic surfaces. More particularly, embodiments of the present invention relate to an apparatus to apply a lightly scratched finish to a metal surface and a method of construction of the apparatus.\n2. Description of the Related Art\nExposed metal is often used as a wall surface, particularly on the exteriors of buildings, such as event centers, coliseums, or museums. It is usually desirable that the surface have a finish that scatters light to avoid harsh glares or strong reflection of the sun. Generally, a lightly scratched metal surface has a soft or matted appearance that is more aesthetically pleasing.\nPrior art techniques to produce a lightly-scratched metallic surface have involved manual approaches as well as machine-controlled processes. With manual approaches, an operator may use a handheld sanding disc or a floor sander to apply scratches to the surface of the metal and must visually judge the level of scratching on the surface. Typically, some areas are more scratched than others, resulting in a surface that looks non-uniform.\nMachine-controlled techniques may involve adapting a rotating sanding disc to a computer numerical controlled (CNC) machine. The sanding disc may be scanned across the metallic surface to apply scratches in a more uniform manner than manual approaches. However, the surface may include subtle lines where the scanning pattern overlaps that have a darker shading than other areas of the surface. Furthermore, the overlapped areas may appear to have large arc-shaped features under certain lighting. As a result, the machine-controlled approach may be less aesthetically pleasing due to subtle imperfections."}
{"text": "The present invention relates to a call management system for a telephone, a database, a telephone and to methods and systems for managing a call.\nWhen a user attempts to make a telephone call, the connection is not always successful. For example, the other party may be engaged on another call; their telephone may be switched off, out of range, or otherwise unavailable; the other party may decline to accept the call; or the other party may not be available to take the call. In domestic situations, when there is a firm intention on the part of the user to contact only the intended other party, such situations are inconvenient but a further attempt at connection is generally made, and successful, after a brief period of time.\nWhen a user is attempting to contact the provider of a service (e.g. a business) the failure to make a connection is especially inconvenient, since the user may wish to obtain the services as soon as possible. As with the above described domestic situation, the service provider may be engaged on another call and unable to take the incoming call, or may be away from the telephone and not be able to answer the call. Even if the service provider has multiple telephone lines, all of the lines may be busy. There are myriad other reasons why the connection may not be possible.\nSuppose that a user's vehicle has broken down and the user attempts to contact a breakdown recovery service. Failure to connect to the recovery service on the first attempt results in the user having to wait by the roadside, making repeated attempts to connect to the recovery service. It is rare that a user who has broken down by the roadside will have ready access to a list of other recovery services in the area. To find alternative services may require the user to leave their vehicle to locate a telephone directory, or make use of a chargeable directory service.\nUltimately, the failure to connect calls results in reduced revenue for the telephone network operator. This is because the network operator does not collect payment from a user for an unconnected call but still incurs overhead costs for the operation of the telephone network infrastructure handling the attempted connection.\nIt is estimated that, on average, for the telephone network as a whole (including fixed line, mobile networks, VoIP etc), over 20% of all voice calls are not connected on the first attempt. With mobile networks, particularly, available signal strengths influence call connections still further and it is estimated that between 30% and 40% of calls including a mobile telephone are not connected on the first attempt. This figure may be still higher for mobile-to-mobile calls.\nThere is a need, therefore, at least to reduce the number of failed connections and to increase the ability of a user to reach a party on the first attempt to establish a connection."}
{"text": "1. Field of the Invention\nThis invention generally relates to a multi-stage or multi-speed planet wheel of a sun and planet gear mechanism having at least two planet wheels, in particular for a bicycle wheel hub.\n2. Background Information\nSuch bicycle wheel hubs with multi-stage planet sun and planet gear mechanisms are disclosed, for example, in United Kingdom Patent No. 21 66 502 B, Federal Republic of Germany Patent No. 814 982, Federal Republic of Germany Laid Open Patent Application No. 41 42 867A1. The individual speeds of the planetary gear are generally selected by momentarily fixing one of the sun wheels, which sun wheels are always engaged with the planet wheels on the stationary hub axle, while the other sun wheels remain free to rotate. An additional shifting capability is frequently created by optionally coupling the planet wheel carrier, generally corresponding to high gears or speeds, or the ring gear, generally corresponding to low gears or speeds, are coupled to the driving toothed wheel of the planetary gear, so that, accordingly, a two-stage planetary gear makes available two high speeds and two low speeds. Frequently, there is also a direct speed, so that a total of five speeds are available. In the case of a three-stage planetary gear, there are accordingly seven speeds.\nIn conventional planetary transmissions, at least three planet wheels are generally used, to achieve a uniform distribution of force. Low manufacturing costs are important, in particular when the planetary gear is used in a mass-produced item such as a bicycle wheel hub. At the same time, care must be taken to achieve sufficient mechanical strength and the least possible wear. These requirements can be met by using sintered moldings. Sintered moldings are manufactured using non-cutting methods, by compressing metal powder at high temperatures, and they tend not to require any finishing operations. If necessary, the material can also be hardened. But there are limitations with regard to the shapes which can be molded, above all because the powder primary materials are pressed inside the mold by the dies which are opposite one another, and are moved toward one another during the molding operation, essentially only in the axial direction. Any shoulders on the mold which project radially inward can have an adverse effect on the compression process. Other shaping methods, e.g. powder metallurgy, injection molding or lost-wax or precision casting, are generally too expensive for mass-produced items such as bicycle wheel hubs.\nThis is also true for potential manufacturing methods which utilize machining operations. On account of the different number of teeth of the planet wheel stages, in connection with the mold gate required for each stage, the economical cutting processes, e.g. cutting with a solid cylindrical milling cutter, are practically out of the question. The manufacture of the multi-stage planet wheel by a numerically-controlled slotting machine is of course possible in itself, but again it is too expensive.\nMaking the planet wheel in several parts tends not to be coherent with the objective of economical manufacture, because it should generally be absolutely guaranteed that all the planet wheels of the planetary transmission have as close to exactly the same orientation of the gear teeth of the two parts relative to one another as possible. Otherwise there can be deformations inside the planetary transmission, since the planet wheels are essentially constantly engaged in the common ring gear and in the common sun wheels. It should generally therefore be expected that when the planet wheels are designed as two-piece parts, some measures have to be taken to guarantee the substantially exact mutual angular orientation of the two parts when they are joined. But that, in turn, is generally not coherent with the objective of the most economical manufacturing process possible."}
{"text": "Field\nThe present disclosure relates to grain bag loaders and more particularly pertains to a new grain storage bag loader apparatus with a highly versatile and functional design providing greater ease of use and efficiency of use.\nDescription of the Prior Art\nThe storage of materials, such as grains, in large bags of lengths up to 500 or more feet on the ground surface in a field is an increasingly popular way to store materials in a protected manner for short and long periods of time. The loading of the materials into, and subsequent unloading of the materials from, the interior of the bags can be difficult due to the weight of the unfilled bags resulting from the significant lengths and heavy walls of the bags. Further, the movement of the grain materials from the harvesting or transporting vehicles to the bag can require a number of specialized apparatus used together which adds to the cost of using such bags and the difficultly of assembling all necessary parts of the system in the field where the bag is to filled with the material. In some cases, using existing technologies, the convenience and cost savings of using the bags for storage may be outweighed by the cost of the equipment necessary for use of the bags as well as the difficulty in preparing and filling the bag."}
{"text": "Field\nThe field relates to stacked integrated device dies and stress isolations features for stacked integrated device dies.\nDescription of the Related Art\nIn various types of packages, two or more integrated device dies can be stacked on top of one another. Stresses can be transmitted between the stacked dies, which can degrade the performance of the package. Accordingly, there remains a continuing need for reducing the transmission of stresses between stacked integrated device dies."}
{"text": "The present invention relates to semiconductor devices and, more specifically, to MOS gate controlled reference (MOS-gated) semiconductor devices formed using a reduced number of masking steps with no critical alignment.\nMOS-gated devices are well-known in the art and include devices such as the MOS-gated devices described in U.S. patent application Ser. No. 08/299,533, filed Sep. 1, 1994 (IR-1113) U.S. Pat. No. 5,795,793, which is incorporated herein by reference. These devices include power MOSFETs, MOS-gated thyristors, insulated gate bipolar transistors (IGBTs), gate turn-off devices and the like.\nThe manufacturing processes for such devices typically include a number of lithographic masking steps which include critical mask alignment steps. Each of these critical alignment steps add manufacturing time and expense as well as provide possible sources of device defects.\nIt is therefore desirable to minimize or eliminate the number of critical alignments necessary as well as reduce the number of masking steps to improve the manufacturing yield and reduce the manufacturing cost."}
{"text": "Generally, a wiring substrate has, on a main surface (front surface) thereof, terminals for connection to semiconductor chips (hereinafter may be referred to as “connection terminals”). In recent years, connection terminals have been provided at high packing density, and the distance (pitch) between adjacent connection terminals has been reduced. Thus, there has been proposed a wiring substrate having an NSMD (non-solder mask defined) structure in which a plurality of connection terminals are arranged in a single opening of a solder resist layer.\nHowever, when a plurality of connection terminals are arranged in a single opening at small pitches, solder applied (through coating) onto the front surface of a connection terminal may flow toward an adjacent connection terminal, and short circuit may occur between these connection terminals. Therefore, there has been proposed a wiring substrate in which an insulating partition wall is provided between adjacent connection terminals so as to prevent flow of solder applied (through coating) on the front surface of a connection terminal toward an adjacent connection terminal (see, for example, Patent Document 1)."}
{"text": "1. Field of the Invention\nThe present invention relates to a protection circuit of a secondary battery, such as a lithium battery, and more particularly to a battery voltage detector circuit of a protection circuit that uses a multiplexer.\n2. Background Art\nIn general, a secondary battery, such as a lithium battery, uses a protection circuit to protect the battery from overcharge or overdischarge. A protection circuit has a voltage detecting unit for detecting the voltage of a battery. Protecting a plurality of batteries connected in series requires a plurality of voltage detecting units corresponding to the individual batteries, thus leading to a larger scale and a higher withstand voltage of the circuit.\nFIG. 2 is a circuit diagram of a battery voltage detector circuit that uses a conventional multiplexer. The battery voltage detector circuit using the conventional multiplexer includes a battery power supply device 11, a battery 11a, switches 21, 22, 23 and 24, flying capacitors 28 and 29, an amplifier 25, an A/D converter 26, and a controller 30. The switch 21 is composed of normally open contacts 21a and 21b, the switch 23 is composed of normally open contacts 23a and 23b, and the switch 24 is composed of normally open contacts 24a and 24b. \nTo detect the voltage of the battery 11a, the switches 21 to 24 are set to an OFF (open) state. In this state, the switch 21 is first set to an ON state and the normally open contacts 21a and 21b are individually set to a closed state. This causes the voltage in the battery 11a to be applied to the flying capacitors 28 and 29, which are connected in series. Thus, electric charges are accumulated in the flying capacitors 28 and 29.\nAfter the switch 21 is held in an ON position for a predetermined period of time, the switch 21 is set to the OFF state and the normally open contacts 21 a and 21b are individually set to the open state. This causes the electric charges corresponding to the voltage of the battery 11a to be accumulated in the flying capacitors 28 and 29.\nThereafter, the switch 22 and the switch 24 are turned on. Turning the switch 22 on causes the connection point of the flying capacitors 28 and 29 to be connected to the ground and fixed to 0V. When the switch 24 is turned on and the normally open contacts 24a and 24b are individually set to the closed state, an inverting input terminal of the amplifier 25 is fixed to the same potential (0V) of an output of the amplifier 25, and the voltage of a non-inverting input terminal is fixed to the ground (0V).\nThereafter, the switch 23 is turned on and the normally open contacts 23a and 23b are individually set to the closed state. This causes the flying capacitors 28 and 29 to be connected to the inverting input terminal and the non-inverting input terminal, respectively, of the amplifier 25. However, the switch 24 is on, so that the voltages of the terminals are fixed and the voltages of the flying capacitors 28 and 29 are not applied to the terminals of the amplifier 25.\nThereafter, the switch 24 is turned off and the fixation of the voltages at the terminals of the amplifier 25 are cleared, thereby causing the voltages accumulated in the flying capacitors 28 and 29 to be applied to the amplifier 25. The voltages of the flying capacitors 28 and 29 are applied in the state wherein the voltages of the input terminals of the amplifier 25 are fixed to 0V. This makes it possible to accurately detect the voltages applied from the flying capacitors 28 and 29 without the risk of the voltages of the input terminals exceeding a permissible range or an output of the amplifier 25 being saturated. In addition, the voltages supplied to the amplifier 25 are unlikely to exceed a permissible range, thus protecting the amplifier 25 from deterioration and damage (refer to, for example, Patent Document 1).\n[Patent Document 1] Japanese Patent Application Laid-Open No. 2001-201548\nHowever, the conventional art has been posing a problem in that, at the instant the switch 22 is turned on, electric charges move to a parasitic capacitance produced between the switch 22 and the flying capacitors 28 and 29 and the holding voltages of the flying capacitors 28 and 29 change, leading to deteriorated voltage detection accuracy."}
{"text": "1. Field of the Invention\nThe present invention relates to a gripping jig for assembling, an assembling device, and a producing method of an assembled body.\n2. Description of the Related Art\nA gripping jig for assembling has been proposed in which a component member of a ceramic luminous body is set in an electrode reception area so as to fix the component member with a pin (see Patent Document 1, for example). In the gripping jig for assembling described in Patent Document 1, two component members fixed to the electrode reception member are prepared; junctions of the component members are opposed to each other; the junctions are abutted to each other while being simultaneously heated for locally melting a cementing material, so that the junctions are joined together by alternately compressing and extending the boundary region between the two junctions.\n[Patent Document 1] PCT Japanese Translation Patent Publication No. 2004-519820"}
{"text": "The present invention relates to combustion turbine engines, and more particularly, to integrating late lean injection into the combustion liner of combustion turbine engines, late lean injection sleeve assemblies, and/or methods of manufacture related thereto.\nMultiple designs exist for staged combustion in combustion turbine engines, but most are complicated assemblies consisting of a plurality of tubing and interfaces. One kind of staged combustion used in combustion turbine engines is late lean injection. In this type of stage combustion, late lean fuel injectors are located downstream of the primary fuel injector. As one of ordinary skill in the art will appreciate, combusting a fuel/air mixture at this downstream location may be used to improve NOx performance. NOx, or oxides of nitrogen, is one of the primary undesirable air polluting emissions produced by combustion turbine engines that burn conventional hydrocarbon fuels. The late lean injection may also be function as an air bypass, which may be used to improve carbon monoxide or CO emissions during “turn down” or low load operation. It will be appreciated that late lean injection systems may provide other operational benefits.\nCurrent late lean injection assemblies are expensive and costly for both new gas turbine units and retrofits of existing units. One of the reasons for this is the complexity of conventional late lean injection systems, particularly those systems associated with the fuel delivery. The many parts associated with these complex systems must be designed to withstand the extreme thermal and mechanical loads of the turbine environment, which significantly increases manufacturing expense. Even so, conventional late lean injection assemblies still have a high risk for fuel leakage into the compressor discharge casing, which can result in auto-ignition and be a safety hazard. In addition, the complexity of conventional systems increases the cost to assembly.\nAs a result, there is a need form improved late lean injection systems, components, and methods of manufacture, particularly those that reduce system complexity, assembly time, and manufacturing cost."}
{"text": "Reactive magnetron sputtering is the sputtering of an elemental target containing one or more elements in the presence of a gas that will react with the target materials to form a compound inside the reactor. One example of reactive magnetron sputtering is adding nitrogen to the reactor while sputtering titanium to form titanium nitride (TiN). Reactive magnetron sputtering for thin film processes has been used in numerous industrial applications. For example, in semiconductor processing, titanium nitride (TiN) has been used as a hardmask for copper dual damascene applications.\nThe reaction mechanism between the sputtered materials and the reactive gas create unique and complex processing stability problems for reactive magnetron sputtering processes. For example, during target sputtering processes, the reactive gas will react with the target surface, chamber kit surface, and at the substrate to form compounds of the deposited film.\nOne of the biggest challenges in reactive magnetron titanium nitride (TiN) sputtering is defect control. A periodic pure titanium (Ti) sputtering is often implemented to mitigate the brittleness of the titanium nitride (TiN) films on the process kit as well as to mitigate titanium nitride (TiN) formation on the target surface. Due to the periodic pure titanium (Ti) sputtering, the target surface and process kit surface can become significantly altered after titanium (Ti) sputtering, which will cause a noticeable shifting in process parameters and film properties.\nAccordingly, the inventors have provided improved methods and apparatus for processing a substrate."}
{"text": "The present invention generally relates to vechicle brakes and more particularly relates to a brake pressure control device for anti-lock pressure control and for automatic brake intervention.\nDE 42 32 311 C1 discloses a hydraulic automotive vehicle brake system with an anti-lock device which, for improvement of the vehicle tracking stability, is equipped with an automatic brake intervention device for driving dynamics control. Special measures are required to be taken both for driving dynamics control and for traction slip control in order to provide the pressure fluid necessary for brake intervention by means of a pump at a sufficient rapidity. One of these measures is, among others, the arrangement of a pressure sensor at the pressure fluid path extending from the brake pressure generator in order to sense the pilot pressure that is respectively generated by the driver in the brake pressure generator.\nOne constructive design variation for arranging a pressure sensor at a special precharging device for an automotive vehicle brake system is shown on page 687 of the ATZ Journal, 96th volume/no. 11. The arrangement of the pressure sensor at the housing of the so-called charging piston unit normally necessitates a complicated cable connection with the control and regulating electronics.\nFurther, DE 195 14 383 A1 describes a brake pressure control device which includes in a housing a plurality of accommodating bores for pressure modulation valves arranged in valve rows. Besides, pressure sensors are accommodated in the housing outside the valve rows, so that the basic contour of the housing and the arrangement of ducts in the block-shaped housing must be adapted to the pressure sensor arrangement chosen. In addition, the cover which closes the housing and the electronics disposed therein must be adapted due to the related increase of the housing dimensions. The drilling operation for the duct system necessitates intricate effort.\nIn view of the above, an object of the present invention is to configure a brake pressure control device of the type initially referred to in such a manner that a pressure sensor arrangement is provided with relatively little effort in terms of construction and manufacture which ensures the simplest possible connection of the pressure sensor arrangement to the existing duct system in the housing, while the original dimensions of the housing are maintained."}
{"text": "The present invention comprises a new and distinct cultivar of Pelargoniumxc3x97hortorum known by the varietal name xe2x80x98PACPALxe2x80x99 (Oglevee No. 5254, Breeder No. P-3364). The new variety was discovered in a selected breeding program in Dresden, Germany. The new variety was first asexually reproduced by cuttings at Olglevee, Ltd. in Connellsville, Pa. The new variety has been trial and field tested at Connellsville and has been found to retain its distinctive characteristics through successive propagations.\nThe cultivar, when grown in a glass greenhouse in Connellsville, Pa. using natural light and temperatures of 60xc2x0 F. night and 68xc2x0 F. day, has a response time of six weeks from a rooted cutting to flowering in a 4 inch pot. The response time was determined on plants grown in soilless media employing constant fertilizer 200-250 parts per million of nitrogen and potassium in full light."}
{"text": "(1) Field of the Invention\nThe present invention relates to an organic nonlinear optical material. More particularly, the present invention relates to a novel organic nonlinear optical material comprising a 4-nitrophenylaminoacetaldehyde acetal, and also to a novel organic nonlinear optical material comprising an acetalized polyvinyl alcohol having the above-mentioned substance as the side chain.\n(2) Description of the Prior Art\nSince an organic nonlinear optical material has a large nonlinearity, it is expected that this material will be used as a base material of an optoelectronic device in the future.\nFor example, urea, 2-methyl-4-nitroaniline, 4-nitro-4'-methylbenzylideneaniline and 3-methyl-4-nitropyridine N-oxide are known as typical instances of the organic nonlinear optical material.\nHowever, when these organic nonlinear optical materials are used as optical elements, the following serious problems should be solved.\n(1) Most of the materials of the 4-nitroaniline series have a subliming property and have no dimensional stability.\n(2) When a semiconductor laser is used as the light source, a nonlinear organic material should be used as an optical waveguide. The above-mentioned organic materials are, however, poor in the moldability and processability for the production of optical wave-guides.\nFurthermore, the above-mentioned organic materials involve a problem concerning electric field orientation."}
{"text": "Incorporating a telephone interconnect call mode of operation onto a trunked system has generally been constrained by the simplex (voice and data signaling) nature of trunked system subscriber units. Similarly, signaling and call processing by fixed-end control equipment is also simplex formatted. While the introduction of duplex voice/RF-capable subscriber units utilizing simplex data signaling with fixed-end control equipment improves the grade of uses interconnect service, their use restricts the amount of dynamic control and call status changes possible while processing a current, duplex voice-only call. Throughout the present specification, the term \"simplex\" is used in the same context as the term \"2-way half-duplex\" may be used, although the meaning is the same: capable of transmitting in either direction, but not in both directions simultaneously."}
{"text": "In electronic instrumentation and signal processing, a time to digital converter (abbreviated TDC) is a device for converting a signal of sporadic or periodic pulses into a digital representation of their time indices."}
{"text": "The invention relates to a thermal management system for a multiple processor computer appliance. More particularly, the present invention enables cool ambient air to absorb heat generated by multiple processors, or other electronic devices, without loss of cooling efficiency while minimizing the amount of space needed for the cooling system.\nIn the rapid development of computers many advancements have been seen in the areas of processor speed, throughput, communications, and fault tolerance. Initially computer systems were standalone devices in which a processor, memory and peripheral devices all communicated through a single bus. Later, in order to improve performance, several processors were interconnected to memory and peripherals using one or more buses. Today computer components have become so small that multiple high-performance processors, memory and disk drives, as well as other equipment, may all be placed on a single baseboard or circuit board. However, computer equipment in general and microprocessors in particular have always been susceptible to damage due to heat buildup.\nSeveral methods have been employed for cooling high-performance microprocessors. A common method of cooling such microprocessors uses a fan heat sink in which an axial fan is attached to the heat sink atop the microprocessor in order to blow air across the heat sink to remove the heat dissipated by the microprocessor. However, these types of fan heat sinks are not that effective and generate their own sets of problems. First, they significantly increased the height of the baseboard or circuit board. Second, they are not effective cooling methods since the ambient air surrounding a microprocessor may be heated by surrounding components on the circuit board.\nAnother method of cooling high-performance microprocessors may be found in U.S. Pat. No. 5,912,802 to Nelson. FIG. 1 is a block diagram of the cooling system 100 as shown in the Nelson patent. The cooling system 100 has a blower 110, a first air duct 120, a first ducted heat sink 130, a second air duct 140, and a second ducted heat sink 150. The blower 110 brings in ambient air from air intake 135. The blower 110 forces air through the first air duct over the first ducted heat sink 130 which is attached to first microprocessor 160. The air moving over the first ducted heat sink 130 absorbs the dissipated heat from the first microprocessor 160 and continues through the second air duct 140. Thereafter, this warmer than ambient air passes through the second ducted heat sink 150 where it further absorbs heat dissipated from the second microprocessor 170.\nIt should be noted in the Nelson cooling system that the air ducts are connected in series and as the air passes over successive heat sinks attached to microprocessors it eventually would reach a temperature where cooling would be ineffective. The second microprocessor\"\"s heat sink receives pre-heated air from the first microprocessor\"\"s heat sink. Hence the second microprocessor would be inadequately cooled. In order for such a serial cooling system to handle several processors or extremely high-performance processors it is necessary to increase the size of blower 110 as well as the size of the air intake 135 and the first and second air ducts 120 and 130.\nTherefore, what is required is a mechanism via which cool ambient air may be brought into contact with each heat exchanger for each microprocessor. This mechanism should not use air supplied to a microprocessor that has been pre-heated by a previous device. Further, the size of the blower as well as the air ducts should be kept to an absolute minimum so as to conserve power and space required by a computer appliance. Still further, the entire computer including processors, disk drives as well as the heat sinks and cooling system must fit within a 1U industry standard height which is 1.75 inches. In addition, this mechanism must be able to be adjusted so that depending upon the heat generated the appropriate volume of cooling air may be passed over the heat sink."}
{"text": "Spark plugs are conventionally furnished in various heat ranges, the heat range being a measure of the thermal capacity of the spark plug. Because different engines develop different amounts of heat in their combustion chambers during their various phases of operation, they are fitted with spark plugs whose heat ranges are most appropriate for the normal operating conditions of the engine. In operation, the electrodes and insulator of the spark plug adjust to a mean temperature corresponding to the engine load. The insulator nose should be at an operating temperature of between 450.degree. C. (842.degree. F.) and 850.degree. C. (1562.degree. F.) for proper operation. Temperatures above 450.degree. C., the self-cleaning temperature are desirable because at those temperatures the deposits on the insulator in the form of soot, carbon residue, and other compounds are burnt off or undergo chemical change. At idle or part load such as during city driving, at intersections, or on downgrades or the like, the insulator nose temperature often drops below 150.degree. C. (302.degree. F.). If the amount of engine lubricating oil entering the combustion chamber due to worn piston rings, intake valves, etc.; if the fuel-air mixture is not too rich; and if this low temperature is not maintained for too long a time; the spark plug will usually still function properly. A temperature of 850.degree. C. at the insulator should not exceeded, because above 900.degree.C. (1652.degree. F.) auto-ignition can occur, and at very high temperatures the electrodes are also subjected to increased attack or destruction by sulphur and other aggressive compounds resulting from the combustion process.\nConventional spark plug characteristics can thus be defined as follows:\nThe lower the heat range index of a spark plug, the higher its resistance to auto-ignition, but the lower its resistance to fouling.\nThe higher the heat range index of the spark plug, the lower its resistance to auto-ignition, but the higher its resistance to fouling.\nBecause of the wide range of operating conditions an engine can encounter, the selection of the proper heat range for a spark plug for a specific engine necessarily is a compromise. Thus, there is always the possibility that the engine will be operated under conditions that are wrong for the heat range of the spark plug installed therein. This possibility can lead to poor performance and engine damage with consequent owner dissatisfaction. To overcome these and other serious disadvantages of prior art spark plugs, it is thus a principal object of this invention to provide a multiple heat range spark plug whose operating temperature is varied automatically such that it operates at an optimum temperature during all phases of engine operation. It is a concurrent object of this invention to provide a heat pipe in the spark plug to maintain it automatically at an optimum operating temperature, with the centerbore of the spark plug insulator itself forming the walls and incorporating the capillary means of the heat pipe such that its most effective heat transfer characteristics are obtained. A further object is to provide a spark plug whose insulator nose operates at a self-cleaning temperature at all times to thereby reduce the danger of combustion misses and misfiring which cause sharp increases in the emission of hydrocarbons, to thus yield advantages with respect to exhaust-gas values in low load ranges while at the same time improving fuel economy."}
{"text": "Technical Field\nThe present invention relates to system output optimization and, more particularly, to key performance indicator optimization in production systems.\nDescription of the Related Art\nIn many production systems (e.g., power plants, chemical plants, etc.), system output is determined by one or more controllable input variables. For example, in thermal power plants, the ingredient configurations of the input fuel (such as, e.g., coal, etc.) can largely affect the output efficiency of the thermal power plants. This output efficiency is usually calculated as the ratio between the electricity generated and the amount of fuel consumed. Therefore, an important goal for such a plant is to determine how to choose an input fuel with desired properties to improve system efficiency.\nBesides such static input there are also dynamic control schemes that tune some variables during the production process. For example, in many production systems, operators need to tune parameters (such as, e.g., voltage, current rotation angle or certain components, etc.). How to determine a combination of such control variables to maximize the output quality (such as, e.g., a Key Performance Indicator (KPI) value) is thus of great interest to production management and optimizations."}
{"text": "In most video surveillance applications, a fundamental activity is the detection of objects of interest, the tracking of those objects through the video sequence, and the determination of inferences from the object's activities. Nearly all commercially available surveillance algorithms that are presently available depend upon motion detection to isolate objects of interest. However, motion detection has its limits, such as the difficulty in distinguishing between different types of moving objects (e.g., cars, people, and trees) and the difficulty in distinguishing between closely aligned objects (e.g., two people walking next to each other).\nTo address these problems, attempts have been made to develop techniques that are specifically designed to detect people in video data. However, even when an algorithm is specifically designed for the detection of people, difficulties persist such as large variations in global and local appearances due to changes in clothing styles and camera angles, the effect of extraneous matter such as backpacks, and the occlusion of parts of a person in a scene by other people or objects."}
{"text": "There are a number of conventional processes for packaging integrated circuit (IC) dice. By way of example, many IC packages utilize a metallic leadframe. One or more dice may be electrically connected to the leadframe by means of bonding wires or solder bumps that have been preformed on the active surface of the die. In general, the dice and portions of the leadframe are encapsulated with a molding compound to protect the delicate electrical components on the active side of the die.\nThe dice within the package generate significant amounts of heat. Different package designs control heat buildup in a variety of ways. By way of example, a heat sink can be attached to one or more of the dice in the package. The heat sink may be exposed on the exterior of the package. That is, a surface of the heat sink is left uncovered by the molding compound. As a result, heat that is generated within the die can efficiently exit the package through the exposed surface.\nThere are a wide variety of other ways to improve the thermal performance of integrated circuit packages. For example, the die can be mounted on a die attach pad or heat spreader that is exposed on the exterior of the package. Thermal pipes can be provided within the package to help transfer heat from the die to the ambient environment. In some designs, the encapsulant does not entirely cover the die but instead leaves a surface of the die exposed. The package size can be increased so that heat can be more easily dissipated by convection along the sides of the package. In some applications, air flow can be directed to cool the package.\nWhile existing arrangements and methods for packaging dice work well, there are continuing efforts to develop improved packaging techniques that provide cost effective approaches for meeting the needs of a variety of different packaging applications."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the fields of proton conductors. More particularly, the present invention relates to solid-state proton exchange membranes.\n2. Description of the Related Art\nFuel cells are electrochemical devices that convert chemical energy directly into electrical energy. Generally, the fuel cell comprises an anode, at which electrons are catalytically removed from the fuel and fed to an external circuit, and protons are catalytically removed from the fuel and fed across a proton exchange membrane to a cathode, where the electrons, protons, and an oxidant are recombined to close the circuit.\nFuel cells require hydrogen, but as is known, free hydrogen is highly reactive. Therefore, there has been interest in fuel cells that extract hydrogen from another fuel at the anode. One fuel of interest is methanol (CH3OH), which can be readily transported as a liquid and has a relatively high energy capacity (e.g., assuming 50% efficiency, 250 mL methanol can deliver about 600 watt-hours of electricity). Ethanol (CH3CH2OH) can also be used; its energy capacity is somewhat lower than that of methanol, but ethanol is far less toxic if imbibed or otherwise absorbed. Methanol or ethanol can be oxidized at the anode by exposure to an oxidant, typically oxygen from air, in the presence of a catalyst, producing CO2, H+, and e−. The skilled artisan will understand this description is not stoichiometric.\nAmong the challenges faced by methanol- or ethanol-fueled fuel cells are the optimization of the properties of the proton exchange membrane. The membrane should efficiently conduct protons from the anode to the cathode to complete the circuit. The membrane also should have a low methanol permeation rate, both to enhance the efficiency of methanol oxidation at the anode and to minimize the reduction of cathode potential by oxidation of methanol at the cathode.\nOne area of particular interest are medium temperature fuel cells, by which is meant those suitable for use at temperatures from about 100° C. to about 250° C. To operate at such a fuel cell temperature, a proton exchange membrane should have a high thermal stability. However, thermal stability is typically inversely correlated with proton conductivity, e.g., a material with a high thermal stability generally has a low proton conductivity."}
{"text": "[Not Applicable ]\n1. Field of the Invention\nThe present invention relates to a golf club. More specifically, the present invention relates to a golf club with a weighting member composed of multiple materials.\n2. Description of the Related Art\nGolf club designs are constantly evolving with the primary purpose to improve a golfer\"\"s performance. While the improvements may address a number of areas, a designer strives to design a more forgiving golf club. Forgiveness in a golf club may be achieved by shifting the center-of-gravity of a golf club to a desirable location, and creating a larger moment of inertia.\nIt is difficult to increase forgiveness in a golf club head composed of a homogeneous or monolithic material, such as stainless steel, since there is a limit on the overall weight of a golf club acceptable to the typical golfer. To overcome this difficulty, designers have resorted to combining different materials (high density and low density) to achieve the desired center-of gravity and large moment of inertia. A very high-density material provides a designer with the greatest freedom in improving the performance of a golf club head since less volume is needed to achieve the proper weighting. The most economical, commercially available material with a very high density is tungsten, which has a density of 19.3 grams per cubic centimeter.\nOne challenge in using heterogeneous materials is the ability to join the materials together in a golf club head. Numerous techniques have been created by the golf industry to join heterogeneou materials in a golf club head. One example is the GREAT BIG BERTHA(copyright) TUNGSTEN-TITANIUM(trademark) irons, developed by the Callaway Golf Company of Carlsbad, Calif., which used a screw to attach a tungsten block to the rear and sole of a titanium iron. Another example is the GREAT BIG BERTHA(copyright) TUNGSTEN-INJECTED(trademark) HAWK EYE (copyright) irons, also developed by the Callaway Golf Company, which feature an internal cavity with tungsten pellets in a solder, as set forth in U.S. Pat. No. 6,210,290, for an Internal Cavity Tungsten Titanium Iron, filed on Jun. 11, 1999. An example of a wood is the GREAT BIG BERTHA(copyright) HAWK EYE(copyright) drivers and fairway woods, also developed by the Callaway Golf Company, which use a tungsten screw in the sole of a titanium club head body. Other techniques use adhesives to join the materials, press fit the materials, braze the materials, or structurally hold one material piece within another material piece using undercuts or pockets.\nFor the most part, these techniques require a precisely machined weighting piece to fit within a precise location on a golf club head. The most economical method is to cast a golf club head body with a cavity for the weighting piece and attaching the weighting piece with a screw. However, casting tolerance are low, and require either machining of the cavity itself, or machining of the weighting piece to fit each cavity. The use of softer materials is undesirable since this creates difficulty in finishing the final product due to smearing of such soft materials during grinding of the golf club head.\nFurther, a co-casting process, where the weighting piece is incorporated in the mold prior to pouring the base metal, is very problematic depending on the materials since the weighting piece is relatively cold when the hot liquid base metal is cast around it causing thermal shock. Also, thermal expansion mismatch of materials is a problem with co-casting of heterogeneous materials. Other problems arise during reshafting, where the golf club head is heated to remove the shaft. Such heating will result in low melting temperature materials (epoxies and solder) to flow, resulting in the possible movement of weighting pieces.\nThe present invention allows for a golf club head to be easily weighted without precisely machined weighting components. The present invention is able to accomplish this by forming a weighting member composed of a multi-component material into the golf club head by applying pressure to the multi-component material while it is in a cavity of the golf club head.\nThe most general aspect of the present invention is a golf club head with a body and a weighing member. The body has a striking plate, a heel end, a toe end and a cavity. The weighting member is composed of a multicomponent material and is disposed in the cavity of the body.\nAnother aspect of the present invention is a cavity back golf club head having a body and a weighting member. The body has a striking plate, a toe end, a heel end and a main rear cavity opposite the striking plate. A top wall, a bottom wall, a heel wall and a toe wall define the main rear cavity. The bottom wall has a second cavity with a predetermined configuration. The weighting member is disposed within the second cavity and occupies the entire cavity. The weighting member is composed of a multi-component material.\nYet another aspect is a method for manufacturing a golf club head. The method includes introducing a multi-component powder/pellet mixture into a cavity on a body of a golf club head, and applying a predetermined amount of force to the multi-component powder/pellet mixture to compact the multi-component powder/pellet mixture within the cavity and form a weighting member. The multi-component powder/pellet mixture may also be heated to a predetermined temperature for hot pressing of the multi-component powder/pellet mixture within. The predetermined temperature is above or closest to the melting temperature of one component of the multi-component powder/pellet mixture.\nThe multi-component powder/pellet mixture may be composed of a high-density component and a binder component. One variation of the multi-component powder/pellet mixture may be composed of tungsten and tin, or alternatively composed of tungsten and a tin bismuth material.\nHaving briefly described the present invention, the above and further objects, features and advantages thereof will be recognized by those skilled in the pertinent art from the following detailed description of the invention when taken in conjunction with the accompanying drawings."}
{"text": "It has become increasingly important to enable home consumers to obtain and use a wealth of information received from outside the home or business and to use such information for viewing or controlling home/business appliances or consumer electronics. In recent times, the Internet and the World Wide Web (WWW) have been important new sources of electronic information. In addition, radio and television broadcasts available via cable or satellite communications are also additional sources of information available for the home or business consumer. Additionally, the conventional telephone network is yet another source of information and a valuable information communication medium.\nAlthough these various forms of information sources and electronic communication media exist, the prior art systems have been unable to effectively combine these information sources and communication media into a readily usable and convenient control device. Because the information sources provide such a wealth of information in a highly divergent set of formats, it has been particularly difficult to combine this divergent information into a convenient and robust control device. Moreover, because each of the various information sources may not be available at all locations, it is necessary to provide a control device that is uniformly useful no matter where it may be deployed.\nVarious forms of prior art devices exist. For example, U.S. Pat. No. 5,710,605 entitled xe2x80x9cRemote Control Unit for Controlling a Television and Video Cassette Recorder with a Display for Allowing a User to Select Between Various Programming Schedulesxe2x80x9d discloses a universal remote control unit for controlling a television and videocassette recorder. The device includes a data input mechanism securable to a telephone line for receiving signals from an external supplier representing television programming schedule items. The device includes a display mechanism for displaying television programming schedule items and television and videocassette operating functions in a hierarchy of textual and scrollable selection lists. The device also includes a user-operable selector mechanism for allowing a user to scroll through the lists and then select television programming schedule items or television and videocassette functions. Although the device provides a wealth of functionality, the user interface is excessively complex and confusing. The user is presented with a myriad of buttons and menu displays, which easily confuse an unsophisticated user. Further, the device provides only a telephone line link for receiving external information via the telephone network.\nFurther, other devices exist in the prior art for receiving information via the vertical blanking interval (VBI) in a cable television transmission. Using the VBI, data can be transmitted across a video cable at the same time as real-time video images are transmitted on the cable. Using this technique, data can be transmitted to a wide audience using the cable television infrastructure. Many conventional set-top boxes are coupled to such cable TV systems and receive data transmitted during the VBI.\nOther prior art systems describe various ways of displaying and manipulating electronic programming guide (EPG) information on a conventional television (TV) monitor. These on-screen television guides necessarily interrupt the standard video signal while the user is manipulating the EPG display on the monitor. Such interruptions affect all those who are viewing the monitor. Further, no privacy can be provided for a user manipulating the EPG information if the information is displayed on the TV monitor.\nThus, a better control device is needed for receiving information from external sources and for using the information for controlling home appliances.\nA hand-held device and system for monitoring and controlling electronic devices is disclosed. In one embodiment, a control device for monitoring and controlling an electronic device, includes: 1) a base unit including a microcontroller and an interface coupled to the microcontroller for receiving external information; and 2) a personal navigator coupled to the base unit via a data link, the personal navigator further including a display component for displaying a dual partition selection list including at least a portion of the external information received from the interface, the dual partition selection list further includes a first selection list and a second selection list, at least a portion of the first selection list and at least a portion of the second selection list being at least partly concurrently displayable on the display component."}
{"text": "Before a catalyst is selected for use in a commercial application, for example hydrocarbon reactions in petroleum refining, a great number of catalysts may be examined for use in the envisioned application. A large number of newly synthesized catalytic compositions may be considered as candidates. It then becomes important to evaluate each of the potential catalysts to determine the formulations that are the most successful in catalyzing the reaction of interest under a given set of reaction conditions.\nTwo key characteristics of a catalyst that are determinative of its success are the activity of that catalyst and the selectivity of the catalyst. The term activity refers to the rate of conversion of reactants by a given amount of catalyst under specified conditions, and the term selectivity refers to the degree to which a given catalyst favors one reaction compared with another possible reaction, see, MCGRAW-HILL CONCISE ENCYCLOPEDIA OF SCIENCE AND TECHNOLOGY, Parker, S. B., Ed. in Chief; McGraw-Hill: New York, 1984; p. 8.\nThe traditional approach to evaluating the activity and selectivity of new catalysts is a sequential one. When using a micro-reactor or pilot plant, each catalyst is independently tested at a set of specified conditions. Upon completion of the test at each of the set of specified conditions, the current catalyst is removed from the micro-reactor or pilot plant and the next catalyst is loaded. The testing is repeated on the freshly loaded catalyst. The process is repeated sequentially for each of the catalyst formulations. Overall, the process of testing all new catalyst formulations is a lengthy process at best.\nCombinatorial chemistry deals mainly with the synthesis of new compounds. For example, U.S. Pat. No. 5,612,002 B1 and U.S. Pat. No. 5,766,556 B1 teach an apparatus and a method for simultaneous synthesis of multiple compounds. Akporiaye, D. E.; Dahl, I. M.; Karlsson, A.; Wendelbo, R. ANGEW CHEM. INT. ED. 1998, 37, 9-611 disclose a combinatorial approach to the hydrothermal synthesis of zeolites, see also WO 98/36826.\nCombinatorial methods present the possibility of substantially increasing the efficiency of catalyst evaluation. Recently, efforts have been made to use combinatorial methods to increase the efficiency and decrease the time necessary for thorough catalyst testing. For example, WO 97/32208-A1 teaches placing different catalysts in a multi-cell holder with the heat absorbed or liberated in each cell being measured to determine the extent of each reaction. Thermal imaging has also been used; see Holzwarth, A.; Schmodt, H.; Maier, W. F. ANGEW. CHEM. INT. ED., 1998, 37, -47, and Bein, T. ANGEW. CHEM. INT. ED., 1999, 3-3. Measuring the heat absorption or liberation and thermal imaging may provide semi-quantitative data regarding activity of the catalyst in question, but they provide no information about selectivity.\nSome attempts to acquire information as to the reaction products in rapid-throughput catalyst testing are described in Senkan, S. M. NATURE, July 1998, 4(23), 3-353, where laser-induced resonance-enhanced multiphoton ionization is used to analyze a gas flow from each of the fixed catalyst sites. Similarly, Cong, P.; Doolen, R. D.; Fan, Q.; Giaquinta, D. M.; Guan, S.; McFarland, E. W.; Poojary, D. M.; Self, K.; Turner, H. W.; Weinberg, W. H. ANGEW CHEM. INT. ED. 1999, 4-8 teach using a probe with concentric tubing for gas delivery/removal and sampling. Only the fixed bed of catalyst being tested is exposed to the reactant stream, with the excess reactants being removed via vacuum. The single fixed bed of catalyst being tested is heated and the gas mixture directly above the catalyst is sampled and sent to a mass spectrometer.\nAttempts have been made to apply combinatorial chemistry to evaluate the activity of catalysts. Some applications have focused on determining the relative activity of catalysts in a library; see Klien, J.; Lehmann, C. W.; Schmidt, H.; Maier, W. F. ANGEW CHEM. INT. ED. 1998, 37, 39-3372; Taylor, S. J.; Morken, J. P. SCIENCE, April 1998, 0(10), 7-270; and WO 99/34206-A1. Some applications have broadened the information sought to include the selectivity of catalysts. WO 99/19724 A1 discloses screening for activities and selectivities of catalyst libraries having addressable test sites by contacting potential catalysts at the test sites with reactant streams forming product plumes. The product plumes are screened by passing a radiation beam of an energy level to promote photoions and photoelectrons which are detected by microelectrode collection. WO 98/07026 A1 discloses miniaturized reactors where the reaction mixture is analyzed during the reaction time using spectroscopic analysis.\nIn order to determine the activity and selectivity of multiple catalysts, arrays of reactors have been designed to simultaneously examine multiple catalysts using the above mentioned analysis techniques. For example, EP 1108467 A2 teaches reactors with removable sections to allow easy introduction of catalyst to the reactor bed. The reactors are sealed using o-rings to allow quick connection of the reactor parts and also provide a reliable seal between the reactor parts and between each reactor and its environment.\nMany reactors available currently are designed for the situation where the feed streams are all of the same phase, for example two feed components that are both gases. Many process technologies and chemistries require higher-pressure gas-phase catalysis, in which a liquid feedstock is vaporized before contacting the catalyst. This may become challenging due to the fact that many seals used for combinatorial arrays have a temperature limitation that is below the bubble point of many reactor inlet compositional mixtures. For example, the long-term temperature limitation on a typical O-ring seal is about 170° C., while the bubble point of C6 to C9 hydrocarbons, for example toluene, at operating pressures of about 300 psig (2172 kPa) to about 450 psig (3220 kPa) are between about 180° C. and about 240° C. at a hydrogen to toluene molar ratio between about 1 and about 3.\nU.S. Pat. No. 5,453,526 B1 teaches a catalytic reactor where liquid media can be continuously introduced, evaporated, and fed to a catalytic reaction. U.S. Pat. No. 3,359,074 teaches a polycondensation system of a single vertically extending column which is transversely partitioned to define, in descending order, a reaction chamber, an evaporator chamber, and a finishing chamber. Two articles, Bej K. S.; Rao, M. S. IND. ENG. CHEM. RES., 1991 30 (8), 1819-1832, and Eliezer K. F.; Bhinde, M.; Houalla, M.; Broderick, D.; Gates, B. C.; Katzer, J. R.; Olson, J. H. IND. ENG. CHEM. FUNDAM., 1977, 16 (3), 380-385 show where additional particles are used to aid in flow distribution before a feed is contacted with a catalyst. What is needed is an evaporator that can be integrated into a process vessel, that accommodates a liquid feed so that the seals will not be compromised during operation of the process vessel, while providing for the feed to be in a vapor phase during reaction.\nHowever, evaporators in general have some inherent problems associated with their operation. One problem associated with evaporators in general is non-uniform mixing of a liquid feed and a gas feed. Non-uniform mixing may occur when both a gas and a liquid are introduced to an evaporator through a common inlet. The dual feed of liquid and gas causes alternating regions of gas entrainment and liquid pulsation being introduced to an evaporator, and therefore regions of low concentration of the vaporized species followed by regions of high concentration of the vaporized species being sent to a reactor bed.\nAnother problem associated with non-uniform vaporization occurs mainly because of a non-uniform flow of liquid into an evaporator. In the case of slower moving flow, a liquid issuing from an orifice into an evaporator can form droplets that detach at a regular periodicity because of the fluid dynamics of the liquid. The periodic formation and detachment of droplets leads to non-uniform vaporization within the evaporator.\nWhat is needed is an evaporator for use in a process vessel that overcomes the problems of non-uniform mixing and non-uniform vaporization associated with evaporators in general."}
{"text": "1. Technical Field\nThis invention relates to seed layers and a process of manufacturing seed layers for casting silicon suitable for use in solar cells or solar modules.\n2. Discussion of Related Art\nKnown steps to prepare seed material for silicon casting use a monocrystalline boule of silicon prepared by the Czochralski (CZ) technique or the float zone (FZ) technique. Both CZ and FZ techniques produce round or circular boules of material. The round boule is cut down to a square or rectangular cross-section and then sliced into pieces for the seed material. The conventional CZ and FZ techniques produce boules of up to a few hundred millimeters with very large boules being about 300-450 millimeters in diameter. Squaring the boule and slicing the seed material results in a square of about 210 millimeters per side from a 300 millimeter diameter boule. The resulting seed material covers less than 50 percent of a bottom of a crucible used in directional solidification.\nWithout excessive care, loosely abutted seed tiles will tend to form highly defective crystal areas both in the fused seed layer and in the grown ingot above. These highly defective areas grow and become problematic for the solar cells produced from the material, typically resulting in a lower solar cell efficiency.\nThere is a need and a desire to produce seed layers corresponding in size to larger crucibles used for directional solidification. There is a need and a desire to produce monolithic seed layers that do not break during normal handling and can be used during several casting cycles. There is a need and a desire to form substantially defect free seed layers at significantly larger scales. There is also a need and a desire to lower the cost per unit volume of cast crystalline silicon by increasing a volume of cast silicon produced."}
{"text": "This invention relates to a method of communicating information in a communication system and is applicable, in particular but not exclusively, to a synchronous optical network employing synchronous digital hierarchies (SDH), or equivalent protocols."}
{"text": "Modern medical practice often has to deal with cardiac arrhythmia caused both by a pathologic heart state such as cardial ischemia, tachycardia and by various pharmaceutical prepartions, for example, strophanthin. In these cases preparations are used possesing antiarrhythmic activity, i.e. preparations curing the crdiac arythmia.\nAntiarrhythmic activity is known to be displayed by quinidine and novocaine amide (procaine amide) (M.D.Mashkovsky, (Lekarstvennye sredstva\" (Medicinal Agents), \"Medicine\" Publishers, Moscow, v.1, 1977, p.p. 370-372). However, these preparations are toxic and posses insufficient antiarrthythmic activity; cure the arrythmia in certain cases only and, furthermore, they cause some side-effects, e.g. drastically reduce the arterial pressure.\nPreparations having pronounced antiarrhythmic activity within a wide range and causing no side-effects are not known in the present-day medical practice.\nAn object of the present invention is to provide a compound, which would possess a pronounced antiarrhythmic activity, cause no side-effects and be less toxic than the prior-art preparations."}
{"text": "In the past, biphenyl dicarboxylic acid dialkyl esters have been produced by preparing corresponding biphenyl dicarboxylic acids by various known production processes and the esterifying the biphenyl dicarboxylic acids with alcohols.\nAs the conventional processes for producing the biphenyl carboxylic acids, there are known a process for causing coupling of a 3-bromo benzoic acid by a treatment with a KOH-methanol aqueous solution in the presence of a Pd/C catalyst (U.S. Pat. No. 2,809,210), a process for preparing a Grignard reagent from 4,4'-dichlorobiphenyl, then causing a reaction of the resultant reagent with CO.sub.2 (U.S. Pat. No. 2,508,022), a process for oxidizing 3,3'-dimethyl biphenyl by using potassium permanganate (Journal of American Chemistry; J. Am. Chem. Soc. 72, 3221 (1950)), a process for iodizing biphenyl, then causing a reaction of the resultant product with carbon monoxide (Japanese Unexamined Patent Publication No. 63-104942), and a process for oxidizing 4,4'-diisopropyl biphenyl with molecular oxygen (Japanese Unexamined Patent Publication No. 63-122645).\nFurther, as another process for producing the above-mentioned biphenyl- dicarboxylic acid, there is known the process for oxidizing 3,3',4,4'-tetramethyl biphenyl with molecular oxygen to synthesize 3,3'-dimethyl biphenyl-4,4'-dicarboxylic acid (Japanese Examined Patent Publication (Kokoku) No. 52-3377).\nNevertheless, these conventional processes for producing biphenyl dicarboxylic acid suffer from many defects, e.g., it is difficult to acquire the materials, extremely complicated and long steps are required, and there are many byproducts. Accordingly, all of the conventional processes in which biphenyl dicarboxylic acid is produced and then biphenyl dicarboxylic acid dialkyl esters are produced therefrom, have a poor practically in industrial terms.\nFurther, it was difficult to selectively manufacture 3,4,3',4'-substituted dimethyl biphenyl dicarboxylic acid dialkyl esters by the conventional processes or by processes provided by combinations of the same."}
{"text": "Priority is claimed on Japanese Patent Application No. 2006-146882, filed May 26, 2006, the contents of which are incorporated herein by reference.\n1. Field of the Invention\nThe present invention relate to a power supply system of a fuel cell vehicle.\n2. Description of the Related Art\nA conventional power supply system of a fuel cell vehicle is known which increases the output voltage of a fuel cell using a first DC-DC converter and input it to a load, wherein if the output electrical power from the fuel cell is insufficient for the requirement by the load, then the deficiency is compensated by inputting additional electrical power to the load from a power storage device via a second DC-DC converter, while if the output electrical power from the fuel cell is excessive for the requirement by the load, then the excessive electrical power is supplied from the fuel cell to the power storage device via the first DC-DC converter and the second DC-DC converter, to charge the power storage device (for example, refer to Japanese Unexamined Patent Application, First Publication No. 2003-208913).\nThe power supply system of a fuel cell vehicle having the above-mentioned construction needs cooperative operations of the two DC-DC converters to supply stable current and voltage to the load; and it has been an object to establish a control method therefore.\nThe present invention was made in view of the above-mentioned circumstances and has an object of providing a power supply system of a fuel cell vehicle which enables a cooperative control of two DC-DC converters to supply stable current and voltage to a load."}
{"text": "1. Field of the Invention\nThe present invention relates to an elastic fixture and an attachment method for mounting a length measuring apparatus on an object to be measured having a different linear expansion coefficient from that of the apparatus. In particular, the present invention relates to an elastic fixture for a length measuring apparatus which is suitably employed in a unit type linear scale fixedly attached, along with a main scale made of glass contained in a frame made of aluminum, to a machine body made predominantly of iron, and which can absorb large thermal stresses caused by the difference in linear expansion coefficients. The present invention also relates to a method of attaching a length measuring apparatus by means of the elastic fixture.\n2. Description of the Related Art\nA unit type linear scale (herein after referred to as a linear scale) of a multi-point fixing type, such as shown in FIG. 1, has been known as a linear scale for feed-back control employed in a relatively large machine tool or industrial machine in which the generation of vibration or shock, or the like is expected.\nIn general, the linear scale (a length measuring apparatus) 10 includes a frame 12 made of aluminum and a detection head 16. The frame 12 contains a main scale made of glass. The detection head 16 contains a detection unit, such as an index scale, which is relatively movable with respect to the frame 12. Both the frame 12 and the defection head 16 are securely attached to a stage sliding unit of a machine body (not shown) through, for example, screws 20.\nWhen the ambient temperature of the linear scale 10 securely attached in the manner described above is changed, thermal stress caused by the difference in linear expansion coefficients is generated between the aluminum frame 12 of the linear scale 10 and the iron machine body. Thus, deflections and distortions are generated in the frame 12 and the glass scale contained within the frame 12. In addition, the screws 20 securely attaching the frame 12 may be loosened, and the relative position with respect to the machine body may be changed. These may result in the deterioration of the measurement accuracy. Therefore, a mechanism for reducing the thermal stress should be provided in the fixing portions of the frame 12.\nIn order to solve the above problems, the present applicant has proposed an example of the attachment mechanism for the frame 12 in Japanese Patent Laid-Open Publication No. 2003-97936. This attachment mechanism is illustrated in FIG. 2 (a perspective view), FIG. 3 (an enlarged front view of a III-portion of FIG. 2), and FIG. 4 (a schematic illustration of the operating state).\nIn this example, an intermediate-portion fixing block 30 for permanent attachment is provided on a central portion of the frame 12 for maintaining the relative alignment with respect to a mating plane 8 of the machine body. In addition, an end-portion fixing block 32 is provided on each of both the ends of the frame 12. The end-portion fixing block 32 includes a parallel plate spring mechanism 34 for reducing the thermal stress connected to the frame 12 via a plate spring mechanism 13 of the frame 12.\nIn the case of employing the fixing block 32 described above, even when the frame 12 is expanded due to the ambient temperature change, the thermal stress is absorbed and reduced (relieved) through the deformations of a plate spring 13A of the plate spring mechanism 13 and plate springs 34A and 34B of the parallel plate spring mechanism 34, as shown in FIG. 4. In addition, since the frame 12 is permanently attached by the intermediate-portion fixing block 30 at the center, the relative alignment with respect to the mating plane 8 is unchanged.\nHowever, the functions, as a plate spring, of the plate spring mechanism and the parallel plate spring mechanism disclosed in Japanese Patent Laid-Open Publication No. 2003-97936 strongly depend on the materials of parts, thicknesses T1 and T2 of the plate spring, the lengths L1 and L2 of the plate spring, and the like. Therefore, the sizes of the parts and the machining accuracy of the parts are restricted, and the parts are difficult to manufacture. As a result, for example, the reduction of the sizes of the end-portion fixing block 32 in the cross-sectional direction in FIG. 3 and the reduction of the parts cost are difficult to achieve.\nIn addition, a unit type linear scale for a manual machine tool is required to be used in relatively severe temperature environments. In order to apply the spring mechanisms to this unit type linear scale for a manual machine tool, the temperature range and the limit of the applicable length should be improved."}
{"text": "1. Field of the Invention\nThe present invention relates to a fixed-focus imaging lens that forms an optical image of a subject on an imaging device, such as a charge coupled device (CCD) and a complementary metal oxide semiconductor (CMOS), and to an imaging apparatus, such as a digital still camera, a cellular phone with a camera, a mobile information terminal (PDA: Personal Digital Assistance), a smartphone, a tablet terminal, and a mobile game machine, on which the imaging lens is mounted to perform photography.\n2. Description of the Related Art\nAs personal computers have become popular in homes, digital still cameras which are capable of inputting image information about photographed scenes, persons, and the like into the personal computers have spread rapidly. Further, a cellular phone, a smartphone, or a tablet terminal in which a camera module for inputting images is installed has been increasing. Such apparatus having an imaging function uses an imaging device, such as a CCD and a CMOS. Recently, because the imaging device has been miniaturized, there has been also a demand to miniaturize the whole of the imaging apparatus and an imaging lens mounted thereon. Further, since the number of pixels included in the imaging device has also been increasing, there has been a demand to enhance the resolution and performance of the imaging lens. For example, there has been a demand for performance corresponding to high resolution of 5 megapixels or higher, and preferably performance corresponding to high resolution of 8 megapixels or higher.\nTo satisfy such demands, it can be considered that the imaging lens is composed of five or six lenses, which are a relatively large number of lenses. For example, Japanese Patent No. 4858648 (Patent Document 1) and Japanese Patent No. 4792605 (Patent Document 2) propose an imaging lens composed of five lenses. The imaging lens disclosed in Patent Documents 1 and 2 substantially consists of, in order from an object side, five lenses of a first lens that has a positive refractive power, a second lens that has a negative refractive power, a third lens that has a positive refractive power, a fourth lens that has a positive refractive power, and a fifth lens that has a negative refractive power."}
{"text": "In certain parts of the world, such as France and Brazil, the \"heart\" or center portion of the stalk of the palmetto palm tree is considered a great delicacy for human consumption. Until now, it has, for the most part, been harvested in the wild, primarily in Brazil, where it is packed fresh in tins for domestic consumption and for export.\nThe palmetto palm comprises a \"meristem\" portion, which is the lowest (approximately) 1.5 meters (4.9.+-. ft.) of the exposed portion of the growing plant, and, above that, a main stalk portion, of which about 2/3+ meter (2.4+ ft.) is useful as a source for palmheart extraction. The meristem portion of the stalk is normally significantly larger in diameter (i.e., it is about 7.6 cm. or 3 inches) than is the upper portion of the exposed stalk, which is about 5.1 cm. or 2 inches in diameter. The two sections are connected by a tapered section of stalk which, in cross section, is roughly in the shape of an inverted, truncated cone. The processing steps for recovering palmheart sections traditionally have included field cutting an approximately 1 meter length of stalk from that portion of growing plant which includes a portion of the upper, main stalk that is about 2/3 meter (2.4 ft.) in length, and about 1/3 meter (1.2 ft.) of the upper part of the meristem portion, and the interposed tapered section, all as one piece. The object in such harvesting and the subsequent processing which takes place is to recover the very inside, or \"heart\" of the stalk from the entire harvested length. That traditionally has been done manually, using a knife to slit longitudinally the bark and the intermediate layer of the upper stalk portion, and by coring the meristem portion.\nThe narrower portion of the cut length of stalk is more or less round in cross section and has an outer bark layer which, while not as hard as some tree bark, is nevertheless very firm and tough. Therefore, the existing practice, as a prerequisite to removing the bark layer, is to steam the field-cut stalks in an autoclave for about 1/2 hour at about 125 degrees C. to soften the outermost or bark layer. More significantly, however, as harvested, the stalks are covered on the outside with fairly long, very sharp thorn-like projections. Even though they tend to lie at a shallow angle to the surface of the stalk rather than at right angles to it, these projections make the stalks difficult to handle and process. The outer bark layer encloses an intermediate layer which is also more or less circular in cross section, tough (but less tough than the bark layer), and, in the upper portion of a stalk whose outer diameter of the bark layer is about 5.1 cm. (2 inches), is about 3.5 cm. (13/8 inches) in diameter. The intermediate layer typically is somewhat larger in the meristem portion, while the innermost or palmheart section is substantially of uniform diameter throughout the upper, tapered and meristem sections. The palmheart itself, which lies within the intermediate layer, is more or less circular in cross section and, in the example given, is about 1.9 cm. (3/4 inch) in diameter. However, the periphery of the intermediate layer is eccentric with respect to the periphery of the outer bark layer. Further, the palmheart itself is eccentric but with respect to the periphery of the intermediate layer and usually with respect to the outer bark layer also. These circumstances are of less moment as to the meristem and tapered portions because coring type techniques may be used to remove the palmheart from those portions. However, these eccentricities, which (in cross section taken through the stalk) normally are not correspondingly positioned, complicate greatly the harvesting of the heart, particularly from the smaller diameter top section.\nTraditionally, removal of the palmheart portions from these upper portions is done by people who use very sharp knives to split lengthwise first the outer bark layer of the stalk which is then torn away from that which underlies it, and then the intermediate layer which is then torn away from the underlying palmheart. The depth of any cut made in either of these layers to remove it from that which is underneath must be made with particular attention so as not to cut through it into the the underlying palmheart portion as well. Further, since the orientation of the eccentricities of the two outer layers is not the same, it is necessary usually to remove the outer and intermediate layers in separate, sequential steps, between which the stalks are reoriented by turning them. By this means, it is assured that the cuts are made through the thickest parts of each such layer, to minimize the possibility of cut-through. This is further complicated by the pronounced increase in the cross sectional \"diameter\" of the stalk moving past the tapered portion to the meristem portion from the upper or main portion of the stalk. (In the context of his disclosure it is to be understood that rarely is any portion of a stalk truly circular in cross section, but rather usually it is merely irregularly round, and thus by its \"diameter\" is meant the average of several mean distances taken across its cross-section). Since there is an annular ring in the region of the tapered portion which is comparatively weak that forms the juncture between the tapered portion and the upper portion of the stalk, usually bending the outer bark layer away from the upper part of the stalk which it surrounds causes it to snap at the juncture and separate from the the remainder of the stalk. Thereafter, the upper portion, now consisting of the intermediate layer-covered palmheart portion is severed from the remainder of the stalk and the palmheart is removed from it by linearly slitting its intermediate layer and stripping that layer away from the underlying palmheart core. It is to the phase of coring the palmheart center component from the remainder of the stalk, which consists of the meristem section and the adjacent tapered portion after the upper portion has been severed therefrom, that this invention is directed.\nConcurrently with and/or sequentially and independently from those operations, the palmheart may be removed from the remainder of the original stalk, which now consists of the bark covered meristem portion and the adjacent bark covered tapered portion, by coring it out, manually, or by using mechanical means such as a tubular knife that is thrust axially to separate and extract the palmheart portion from it.\nObviously, all of this is very labor intensive, and since the recovered palmheart cannot be preserved for an appreciably length of time, must be done at or very near the place of harvesting. For these reasons, the price for this commodity traditionally has been very high. Recently, there have been efforts to raise palmettos for palmheart harvesting on plantations, to make the growing and harvesting of them easier and less expensive and to facilitate the rapid transport of cut stalks to processing locations. This also produces products of higher quality and better uniformity than can be obtained with cuttings from growth in the wild. The less wide range in \"diameter\" of plantation stalks, whose ages are more easily unified than is possible with wild stock, has led to interest in automating selected phases of the process, since this also can have the additional advantages of improving the quality of the product and making processing less expensive.\nAccordingly, it is an object of this invention to provide means to recover palmhearts from the meristem and tapered portions of palm stalks.\nAnother object of this invention is to provide such means which is mechanized.\nStill another object of this invention is to provide means for satisfying one or more of the foregoing objectives which is adapted to accommodate substantial diameter variations between and within individual palm stalks as they are being processed."}
{"text": "Processing techniques and facilities which enable widespread use of molded thermoplastic composite components at production rates and production costs and that allow significant weight savings scenarios may be desirable in some applications. The capability to rapidly heat, consolidate and cool in a controlled manner may be required for high production rates of composite components. Current processing techniques include the use of heated dies, and therefore, may not allow for the optimum controlled cool-down which may be required for optimum fabrication. Furthermore, current processing techniques may have limitations in forming the desired components since such techniques have limitations in the capability to hold the dimensions of the component accurately or maintain the composite in a fully consolidated state and may not optimize performance of the current resin systems.\nSuperplastic forming and hot forming methods for fabricating aluminum and to some extent magnesium components may be hampered by the inability to effectively integrate the superplastic forming process with the heat treatment requirements. The savings produced by the excellent formability at SPF temperatures may be nullified by the loss of dimensional control due to the need to solution-treat and quench the component after superplastic forming to produce competitive strength characteristics.\nThe lower strength of non-heat treatable alloys may be a significant contributing factor mainly as to why there has not been widespread implementation of the SPF of aluminum components in the aerospace industry. Moreover, the long cycles and low strength of characteristic of the current process may be deterrents to using the SPF of aluminum and magnesium in the automotive industry.\nTherefore, an induction forming process is needed which uses induction heating to allow rapid heating and cooling of aluminum and magnesium alloys in combination with the application of tool pressure to form and/or mold near net shaped parts at elevated temperatures with subsequent incorporation of an in-tool quenching step."}
{"text": "This invention relates generally to cosmetic stick holder and applicator devices such as lipstick cases, and more particularly to devices of this type which employ a flexible strap operator to actuate a product-carrying movable cup.\nIn the past a large number of cosmetic or lipstick holders and applicators have been proposed and produced. One class of prior device of this type employed a flexible push-pull strip or band which was fastened at one end to the cup and adjacent its other end to a finger piece which extended outwardly through a slot in the wall of the casing. The strip traveled along a guide having a 180.degree. bend, such that sliding movement of the finger piece in a direction away from the open end of the casing caused advancement of the cup and stick product whereby the latter would protrude through the open casing end. In several of these prior constructions, the free end of the actuator strip was made sufficiently long so as to be capable of extending across the open end of the casing, the strip being also wide enough to constitute a closure therefor when the cup was moved to its retracted position.\nThere were a number of disadvantages and drawbacks in these prior devices. In virtually all constructions heretofore proposed, the cases were either wholly or at least in part, of square or rectangular cross section. This was due to the fact that the actuator strips employed therein were of appreciable width whereby their side edges extended along widely spaced, longitudinal portions of the casing. Such rectangular dispensers were difficult to manipulate and use, since they could not be easily twirled in the manner of a cylinder, in the hand of the user during application. It is well recognized that such rolling or twirling movement is desirable to facilitate the application of lipstick wax, for example.\nIn addition, these prior devices were frequently complex and difficult to assemble. Where the prior constructions employed mating casing halves, the strip and cup had to be applied to one casing half during the assembly and held in a precise, given position while the other casing half was applied. Frequently the strip and cup would undergo a slight shift in position during this step, causing difficulty in applying the second casing half. As a result the assembly time was often excessive. Finally, the prior devices were large and cumbersome in a relative sense, and were not in keeping with the aesthetic requirements of this type of cosmetic article."}
{"text": "Radar systems are systems that use radio frequency (RF) waves to derive information about targets, if any, in a region of interest. An RF pulse may be transmitted into the region of interest and, if a target exists in the path of the pulse, reflected energy will be returned to the radar. The radar may then analyze the reflected signal (and possibly other reflected signals) to determine information about the target. Radar systems are used in a wide variety of different applications involving a variety of different target types. For example, as air traffic control radar may detect and track aircraft about an airport to support air traffic controllers. Likewise, a weather radar may be used to detect and track potentially harmful weather formations in a region of interest. An early warning radar (EWR) may be used to, for example, detect and track incoming enemy missiles and/or aircraft as early as possible to allow counter measures to be deployed before arrival thereof.\nDepending on application, it may be very important to achieve accurate and reliable results in a radar system. For example, a missed detection or inaccurate target position/velocity in an EWR can have devastating results. Therefore, radar systems are often required to meet tight operational specifications to ensure reliability and accuracy. If a particular radar system does not meet, the specifications, changes may be required to the system design to bring it within specification."}
{"text": "When camping with a trailer or a van it is desirable in good weather to prepare meals outdoors. Where there are no picnic tables available it is usual to set up a collapsible table, such as a card table, stowed in the camper or van, on which a stove and cooking utensils are rested. It is, of course, time consuming and inconvenient to gather all the materials together for this purpose and to return them to their places in the camper or van after use.\nIt is an object of the present invention to provide a kitchen table which is collapsible as a storage unit for cooking utensils and the like, and which is easily set up for use with the utensils being at hand."}
{"text": "Power consumption of electrical devices in the office or home can be significant. In many instances, devices such as computers remain on for the entire work day, and in some instances, continuously for twenty four hours everyday. In addition to the computers themselves, attached computer peripherals such as display devices and printers also consume significant amounts of power and often remain powered on regardless of the operating state of the computers. Additional example devices requiring power in the home or office include copiers, fax machines, environmental systems, and lighting systems. As energy costs increase, it is desirable to reduce unnecessary power consumption in such devices.\nAs a result, improved methods and apparatuses for power management of electrical devices are needed."}
{"text": "The invention relates to a double layered paper making forming fabric with a coarse structured running side and a fine structured paper side.\nSince long time the paper makers request that in the manufacturing of paper making forming fabrics the paper side, i.e. the paper making surface of the paper making forming fabric, should be provided with an as small meshed structure as possible to obtain an optimal fiber depository at a high as possible retention and a complete absence of marking, while the running side of the paper making forming fabric should be provided with coarse meshes to achieve a good dewatering of the pulp, whereby furthermore the susceptibility to pollution should be reduced and the cleaning possibilities should be improved. Additionally the running side should be sufficiently wear resistant by respective built-in-materials, so that at least the life time usual nowadays for double layered forming fabrics is achieved.\nFor the attainment of the above quality features of double layered paper making forming fabrics considerable progresses have been made in the last years, but no fabrics could have been manufactured which are adequate in their entirety to the modern, quick and efficient but also more sensitive printing processes to that extent that they can deliver papers of highest quality. Even the so-called composite forming fabrics, comprising two complete fabrics connected to one another could not completely fulfill as yet the requirements of the requested paper qualities.\nIn this context attention is drawn for instance to U.S. Pat. No. 4,112,982 disclosing two-layered forming fabrics which are characterized by transverse or cross floatings on the paper side and a coarse structure of the running side by means of relatively thick cross threads of a relative low number. The practice has shown in this connection that only cross threads of a similar thickness can be interwoven with the fine longitudinal threads of the known double layered fabric on the running side thereof and that by increasing the diameter of the cross threads the weft runner quality is lost more and more. That means, the distance between the wear plane provided by the cross thread floatings and the plane of the longitudinal thread knuckles is getting less and less. That is the reason why the longitudinal threads are exposed too early to the wearing so that a quick loss of the stability in the longitudinal direction will occur leading finally to a tearing in transverse direction of the fabric and a reduced running time."}
{"text": "1. Field of the Invention\nThis invention relates to a structure of a sun roof of an automotive vehicle and a method of the assembly thereof.\n2. Description of the Related Art\nThere is an increasing demand of an automobile with an opening formed in a roof panel and for the roof opening structured so as to be opened and closed with a movable roof member.\nThe opening to be formed in the roof panel of the automobile may be structured in such a manner that, as disclosed in Japanese Patent Laid-open Publication (kokai) No. 2-306,821, a pair of left-hand and right-hand front roof opening sections are provided in the front portion of the roof panel. Each of the left-hand and right-hand roof opening sections is open to its entire length at its transversely outer portions, that is, at its left-hand and right-hand outer side ends, respectively. It is closed at its front and rear edge portions and a transversely inner edge portion extending longitudinally (from front to back) so as to connect its front edge portion to its rear edge portion. The left-hand and right-hand front roof opening sections are configured each in a fashion such that a front edge portion of the roof opening section is composed of a front header, a rear edge portion thereof is composed of the front edge portion of a rear roof panel section, and a transversely inner edge .portion extending longitudinally is composed of a connecting middle section that connects the front header with the rear roof panel section. This type of the roof opening may be called a T-bar roof, and it can provide the .passenger with an open feeling that resembles the feeling given by an open car because the left-hand and right-hand outer side ends of the roof panel can be opened.\nThe left-hand and right-hand front roof opening sections are opened or closed with detachable left-hand and right-hand roof sections, respectively, in many occasions. They are also designed to be opened or closed simultaneously with one sheet of a movable roof member that has a size substantially large enough to cover a pair of the left-hand and right-hand front roof opening sections and that is structured so as to be movable in a longitudinal direction of the vehicle, i.e. in forward and rearward directions.\nIn the latter case, that is, when .the left-hand and right-hand front roof opening sections are to be opened or closed simultaneously with one sheet of the movable roof member, the movable roof member should be guided in the forward or rearward direction. At this end, a base frame may be fixed to a body of the vehicle and a guide mechanism may be disposed in a position between the base frame and the movable roof member.\nAlthough the base frame may assume a variety of forms, it is desired from the point of view of ensuring a strength of mounting and the like that the base frame comprises a front base frame section extending over the substantially entire length along and parallel to the front header, a rear base frame section extending over the substantially entire length along and parallel to the rear edge portion of the rear roof panel section, and a middle base frame section extending longitudinally and connecting a middle portion of the front base frame section with a middle portion of the rear base frame section. The arrangement for this structure is preferred in terms of ensuring a wide area for mounting the base frame to the body of the vehicle and ensuring a large number of positions for fixing the base frame to the body thereof because the base frame can be fixed tightly to the body. For this structure, the guide mechanism for the movable roof member may be configured by taking advantage of the middle portion of the base frame.\nOn the other hand, when the guide mechanism for the movable roof member is to be disposed on the rear side of the body of the vehicle, the base frame may be structured in such a shape as having a pair of left-hand and right-hand guide sections, each extending from a transversely outer end portion of the rear base frame section toward the rear.\nFrom the description as have been described hereinabove, the base frame may be configured in such a manner that it comprises the front base frame section, the rear base frame section, and the middle base frame section, or the front base frame section, the rear base frame section, the middle base frame section, and the left-hand and right-hand guide sections.\nThe automotive vehicle having a pair of left-hand and right-hand roof opening sections is provided with no roof side rail at their transversely outer edge portions. This leads to the fact that the transversely outer end portions of the front header, that is, an upper end portion of a front pillar, may deviate from its normal position.\nOn the other hand, since the front base frame section of the base frame is of such a cantilever type as being free at the transversely outer end portions, the transversely outer end portions of the front base frame section are likely to deviate from their normal positions.\nLikewise, the left-hand and right-hand guide sections of the base frame are likely to deviate from their normal positions.\nAs have been described hereinabove, the deviation in the positions is likely to occur on the side of the body of the vehicle as well as on the side of the base frame, so that the mere fixing of the base frame to the body may deviate from their normal positions, thereby failing to guide the movable roof member smoothly. Further, if the movable roof member would be closed in a deviated state, a problem with sealing may arise."}
{"text": "There is an increasing number of consumers who are consuming media, including video content (e.g., movies and TV shows) via streaming media services. Often times, consumers wish to receive information about media content they are watching. Current methods of receiving information about media content advancing through media content might result in a less than ideal user experience.\nFor simplicity and clarity of illustration, the drawing figures illustrate the general manner of construction, and descriptions and details of well-known features and techniques might be omitted to avoid unnecessarily obscuring the present disclosure. Additionally, elements in the drawing figures are not necessarily drawn to scale. For example, the dimensions of some of the elements in the figures might be exaggerated relative to other elements to help improve understanding of embodiments of the present disclosure. The same reference numerals in different figures denote the same elements.\nThe terms “first,” “second,” “third,” “fourth,” and the like in the description and in the claims, if any, are used for distinguishing between similar elements and not necessarily for describing a particular sequential or chronological order. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments described herein are, for example, capable of operation in sequences other than those illustrated or otherwise described herein. Furthermore, the terms “include,” and “have,” and any variations thereof, are intended to cover a non-exclusive inclusion, such that a process, method, system, article, device, or apparatus that comprises a list of elements is not necessarily limited to those elements, but might include other elements not expressly listed or inherent to such process, method, system, article, device, or apparatus.\nThe terms “left,” “right,” “front,” “back,” “top,” “bottom,” “over,” “under,” and the like in the description and in the claims, if any, are used for descriptive purposes and not necessarily for describing permanent relative positions. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the apparatus, methods, and/or articles of manufacture described herein are, for example, capable of operation in other orientations than those illustrated or otherwise described herein.\nThe terms “couple,” “coupled,” “couples,” “coupling,” and the like should be broadly understood and refer to connecting two or more elements mechanically and/or otherwise. Two or more electrical elements can be electrically coupled together, but not be mechanically or otherwise coupled together. Coupling can be for any length of time, e.g., permanent or semi-permanent or only for an instant. “Electrical coupling” and the like should be broadly understood and include electrical coupling of all types. The absence of the word “removably,” “removable,” and the like near the word “coupled,” and the like does not mean that the coupling, etc. in question is or is not removable.\nAs defined herein, two or more elements are “integral” if they are comprised of the same piece of material. As defined herein, two or more elements are “non-integral” if each is comprised of a different piece of material.\nAs defined herein, “approximately” can, in some embodiments, mean within plus or minus ten percent of the stated value. In other embodiments, “approximately” can mean within plus or minus five percent of the stated value. In further embodiments, “approximately” can mean within plus or minus three percent of the stated value. In yet other embodiments, “approximately” can mean within plus or minus one percent of the stated value."}
{"text": "The present invention relates to an replacement process upon a recording failure in an optical disc having a plurality of layers, an optical disc apparatus and an information recording method for recording information onto the optical disc.\nFor a DVD-RAM, a Blu-ray Disc (hereinafter, referred to as BD) and an HD DVD Disc (hereinafter, referred to as HD) which can perform recording, their specifications describe an replacement process to be performed upon a recording failure in a data region. After information is recorded and the recorded information is reproduced, the reproduction signal is evaluated. If the reproduced signal is defective, the “replacement process” is executed. That is, upon detection of a recording failure, the replacement process records the same information as the recorded information in a predetermined region on the disc. The predetermined region will be referred to as a spare area and the reproduction signal quality evaluation operation will be referred to as verify. The spare area is arranged in the inner circumference and the outer circumference of the DVD-RAM, the BD, and the HD. The replacement process and improvement of its reliability are detailed in JP-A-2006-85797 and JP-A-2006-185509.\nOn the other hand, recently, in order to increase the recording capacity, various types of recording discs have multiple layers. Currently, two-layer recording media are normalized in the DVD-R/RW, +RW, BD. Unlike the single-layer disc, the multi-layer disc should consider affect to the other layer upon recording. The multi-layer recording technique is detailed in JP-A-2000-311346.\nSince recording into the spare area is for protection upon a recording failure in the data region, it is necessary to increase the reliability of the recording operation as compared to recording operation in the normal data region.\nOn the other hand, when the optical disc has multiple layers, the following occur as affects from other layers.\n[1] Leak-in of a stray light signal is generated by irregularities of inter-layer distance in the disc plane.\n[2] Leak-in of a stray light signal and an information signal is generated by a pseudo-focal point and focal point-shifted spot in the layer other than the recording or the reproduction layer.\n[3] Recording irregularities are generated by a difference in transmittance in other recorded and unrecorded layers.\nThe leak-in of the stray light by [1] or [2] lowers the signal S/N, which in turn lowers the servo stability and lowers the reliability of the recording and reproduction operation. Furthermore, when the inter-layer thickness is made constant so as to reduce the multi-layered disc manufacturing cost, the leak-in of the pseudo-signal of [2] is generated. This causes fluctuations by leak-in of a servo error signal and lowering of the signal quality by leak-in of a reproduction signal. In such a condition, an error is caused in the reproduction quality evaluation of a recording signal upon verify and it becomes impossible to guarantee a correct verify operation. As for [3], since the recording film state is different in the mark and the space portion of the data recording portion, the transmittance differs and the recording may not be performed uniformly. These affects become greater as the number of layers in the disc increases and the inter-layer gap decreases.\nThe aforementioned JP-A-2006-85797 and JP-A-2006-185509 disclose a technique of data rearrangement and modification of a recording position so as to improve the reliability of the replacement process but do not consider the aforementioned problems in the multi-layer disc. Moreover, as a method for improving the reliability among the aforementioned objects, JP-A-2000-311346 discloses a correction method for recording power learning in the recorded and unrecorded portions of the multi-layer disc. According to the method, it is possible to improve the reliability of the recording learning, or recording parameter adjusting, in the multi-layer disc. However, no consideration is taken on the improvement of recording reliability for the recording operation, especially in the recording in the spare area."}
{"text": "Electricity and other forms of energy are billed in accordance with energy consumption that may vary with the size of the customer and the time of day. Since the peak amount of energy consumed by a customer over a given time period (so-called \"demand\") determines the size of the service required, e.g. size of conductors, transformers, peak generating capacity, etc., many utilities measure this peak consumption or demand to determine the rate to be charged the customer for all electricity consumed over a given period.\nTo determine the amount of energy being consumed by each customer during successive demand intervals, utility companies have located, at customer sites, electricity meters, such as watthour meters containing a demand register, which must periodically be read either by the customer or by a representative of the utility, to accumulate billing data (e.g. peak demand and total energy consumed). In addition, to evaluate customer energy demand so as to assess the capability of the utility's equipment to satisfy demand, or to justify rate modification, energy demand occasionally must be surveyed over an extended period of time, e.g. 18 months.\nElectrical demand is typically measured by the use of a pulse initiator which utilizes a photo-optical detector to detect the rotation of the eddy-disk of a watt-hour meter and produces a series of pulses whose frequency is directly related to the instantaneous power (demand) being delivered to the customer.\nTypically, a demand register accumulates these pulses over a preselected interval, e.g. 15 minutes, to give an indication of the peak demand over the interval. This peak demand data is then stored in a memory device for subsequent readout or display.\nSince there will often be large numbers of meters registering demand throughout a utility system there exists a need to monitor electrical demand simultaneously at a large number of individual customer sites, and to process the data at a central location. This is commonly done by means of a mainframe computer operated by the utility company. This information is used by the utility to assess demand parameters such as peak demand and seasonal as well as daily demand variations. Accordingly, systems have been developed for polling customer meters to extract and transmit energy demand data over commercial telephone lines to the central computer for accumulation and processing.\nBecause the investment required to provide automatic polling of remote meter registers at the customer site is substantial, as an alternative, automatic meter data recorders have been developed that are set up at customer sites to be monitored. The data recorders commonly employ magnetic tape storage techniques to record pulses generated by the pulse initiator of the meter that represent electrical demand during successive intervals. The magnetic tape is then carried to the utility for processing.\nMagnetic tape recorders of the type employed for this purpose are relatively complex, require substantial battery power to operate in the event of a power failure and may fail to work properly in environments exposed to temperature extremes. Recently, due in part to availability of inexpensive solid state memory devices, portable meter reading devices having solid state memories have been developed for accumulating demand data stored in a meter. The reading device is then either brought to the central computer and the data transmitted directly to the computer, or the data is transmitted over telephone lines to the central computer. There is a tendency, however, for errors to occur in the acquisition and transmission of data. These errors are caused by among other things failures of the solid state memory, misdetection of pulses generated by the electricity meter and, most commonly, electrical noise contamination. Data acquisition errors are particularly troublesome because they affect, among other things, the accuracy of customer billing. For this reason there currently exists a significant need to verify energy demand data accumulated from customer meters to ensure that the data is accurate. Also, since peak demand data must be correlated with the time of day, day-of-week, season, etc. for proper billing purposes, it is necessary that such a data acquisition device be capable of operation despite power outages occurring on the power line.\nIn systems which employ battery back-up power supplies for operation in the event of a power failure, there exists a need for determining when the battery has grown weak and requires replacement. Problems can occur in accumulating and maintaining verifiable energy demand data if the battery back-up system fails to provide a reliable alternative source of power during power outages. If the back-up battery has become marginal or failed (especially during periods of frequent outages), the recorded customer demand data may be erroneous or suspect. Yet, the utility has no way of determining when the back-up battery failed so as to estimate the beginning of the suspect data. Prior to the present invention, there has been no known way to accurately monitor the performance of battery back-up systems short of regular planned replacement at periodic intervals in accordance with known statistical replacement techniques, or periodic on-site testing of the batteries at the customers' facilities.\nOther problems have been encountered in conventional remote meter data recorders which transmit energy demand data over commercial telephone lines to a central computer for processing. Because of this preexisting communications link for meter data collection, it would be economically advantageous to the use the same equipment for energy management functions, such as where customer loads are interrupted during periods of peak demand. However, some prior art data recorders are configured to transmit the energy demand data to the central computer only at predetermined times. These predetermined times are usually selected when a customer's telephone line is not likely to be in use, e.g. during the late night hours. This type of prior art data recorder can receive command information from the central computer over the telephone line only after the recorder has initiated the communication by dialing up the central computer to download its data. Systems which can receive commands only during recorder-initiated communication modes are not receptive to receipt of commands at any time other than when communication has been initiated by the recorder, and are therefore unable to provide energy management functions.\nStill other problems have been encountered in data recorder systems which are receptive to commands only after initiation of a communication by the meter for the purpose of downloading data. In particular, there exists a need for detection of error conditions such as power outages, low battery condition, or a malfunction in the meter. Conventional meters and recorders which are unable to provide indications of error conditions have proven unsatisfactory to electrical utilities which need to be able to monitor and correct error conditions rapidly and to perform peak load energy management functions.\nOne object of the invention, therefore, is to provide a method of and system for verifying data accumulated from remote electricity meters during successive demand intervals and to confirm that the data measured represents data upon which a utility company or other controlling institution can rely. Another object is to provide a relatively inexpensive and reliable device for making energy demand measurements and verifying their reliability at customer meter sites for the purpose of obtaining demand survey and customer billing data. A further object is to provide an electricity demand data acquisition device which operates reliably even during a power outage.\nAnother object of the present invention is to provide an electricity demand acquisition device which initiates a communication to a central computer in the event of an error condition such as a power outage, low battery condition, or malfunction in the metering device.\nIt is another object of the present invention to provide an electricity demand data acquisition device which is operative to override the conventional data downloading operational sequence (which occurs at a predetermined time) and receive communications initiated by the utility central computer.\nIt is another object of the present invention to provide an electricity demand data acquisition device which may be configured for operation over a wide range of utility customer demands by scaling the KYZ pulse input.\nIt is another object of the present invention to provide an electricity demand data acquisition device which maintains an indication of the time of day and date of the occurrence of a low battery condition so that an estimation of the reliability of the data stored in the device can be made."}
{"text": "The importance of utilizing natural clean renewable energy has emerged as critical in order to combat global warming. The major sources of natural clean energy are the sun, the ground, wind and water.\nThe use of natural energy to harness power from different sources is limited for different reasons. The use of the sun and wind energy is mostly limited by the Inconsistent availability of the energy source, depending upon weather, seasonal changes, and day and night cycles. The use of ground energy is limited by geographical location and by difficulties in drilling a few miles below the surface of the ground. Water is the most abundant source of clean energy, but the use of energy derived from water is limited by conventional technology that requires dam construction, high water flow, usually exceeding a few meters per second, and the complexity of corresponding energy harnessing devices.\nKinetic energy sources of water movement can be divided largely into three categories, (1) horizontal movement resulting from height differences between two locations in a river, or by tidal flow in the ocean, (2) vertical movement of water in a human built dam or fall, and (3) oscillatory wave movement originated from a combination of the horizontal and vertical movement of water, found mostly in the ocean.\nCurrently, the use of horizontal hydro-kinetic energy in a traditional water mill or propeller-type generator requires a minimum water flow of a few meters per second and is greatly restricted by geographical location. The method also lacks the ability to harness massive power. Several methods have been devised to harness massive energy, in mid-scale to large-scale power generation, based upon horizontal hydro-kinetic energy, including placing a turbine above a river, under a bridge. The method blocks natural water flow, similar to a dam, with the possibility of an occasional increased risk of flooding. The ability to harness horizontal hydro-kinetic energy with no or minimal water flow hindrance is a highly desirable goal. The present invention enables the harnessing of such energy with no or minimal water flow hindrance.\nCurrently, use of the vertical movement of water to harness hydro-kinetic energy is largely limited by the use of an existing natural fall, like Niagara, or by building a dam to hold water above a dam area to create height for harnessing hydro-kinetic energy. Building a dam has mixed advantages and disadvantages. Disadvantages include (1) changes, or even destruction, of an ecosystem, and (2) a requirement for a large scale relocation of people in the affected region. Harnessing vertical hydro-kinetic energy is highly desirable, and the present invention enables harnessing vertical hydro-kinetic energy by recreating height differences without the necessity for building a dam.\nCurrent uses of oscillatory hydro-kinetic energy in several applications are summarized in http://www1.eere.energy.gov/windandhydro/hydrokinetic/default.aspx. These applications still require improvements in efficiency and ease of installation. The present invention provides a simple, highly effective and efficient way of harnessing oscillatory hydro-kinetic energy."}
{"text": "In hydraulic power assist steering systems, a hydraulic pump provides hydraulic fluid for assisting steering action. The pump is directly driven by an engine or an electric motor in the case of an electro-hydraulic power assist system. In these systems, the steering system can be turned to such a degree that the steering rack reaches its travel limits, also referred to in the art as its end stops. If an operator holds the steering system \"against\" the end stops, especially with force, considerable power can be consumed by the pump, while causing no change in the vehicle steering.\nIn systems employing engine driven hydraulic pumps, relief valves have been used to divert hydraulic flow to prevent overloading the hydraulic pump. However, due to their considerable expense and complexity, coupled with their lack of complete effectiveness reducing power loss, other solutions would appear to be desirable.\nIt is therefore a recognized desire among those skilled in the art to provide a hydraulic power assist steering control system in which efficiency is significantly improved by limiting or reducing hydraulic pumping losses otherwise incurred when the system is operated near or against the steering end stops."}
{"text": "This invention relates to liquid purification and separation, for example to the dewatering of a slurry or the like from sewage treatment plants, where a flocculation compound is mixed with a slurry. The solids of the mixture are filtered from the liquids of the mixture, and a relatively dry sludge cake is discharged from the system.\nIt is desirable to dewater slurries in various manufacturing processes so as to recapture and use the solids and to dispose of the water in an uncontaminated condition. For example, in the disposal of sewage, a sludge slurry is removed from settling tanks of sewage treatment plants which contains both solids and liquid, with a substantial amount of solids suspended in the liquid. One common procedure for treating sludge is to add a flocculation compound as a filter aid to the sludge to free the trapped liquid and filter the solids from the liquid. The flocculants are added to the slurry for increasing the filterability of the slurry and for coagulating the solids, therefore reducing the tendency of the small solids from passing with the liquid through the filter medium. Even after flocculation of the sludge mixture some of the solids are so small and the quantity of liquid so great that the smaller solids tend to move from the mixture with the liquid through the filter medium if the liquid is urged too vigorously from the mixture. This condition requires that the sludge be carefully handled and gently dewatered.\nIn the past, several sludge dewatering systems have been developed which comprise pairs of moving conveyor belts usually fabricated from woven textile material that function as filters. The prior art systems deposit coagulated sludge on one moving belt, and as the sludge moves with the first moving belt the second moving belt is moved down into contact with the sludge so as to squeeze and therefore dewater the sludge. Examples of such prior art systems are disclosed in U.S. Pat. Nos. 3,800,952, 3,601,039, 3,699,881, and 3,896,030.\nIn a sewage treating process it is desirable to process as much sludge as possible in a short period of time, to extract as much liquid as possible from the sludge during the process to produce a dry sludge cake, and to capture a high percentage of the solids from the original slurry. These three objectives are in conflict with each other since more water can be extracted from sludge if the sludge is given more time to drain. Although the dewatering capacity of a dewatering system can be increased by merely speeding up the system, the increase in speed also usually causes a lower percent of solids to be filtered from the liquid and the sludge is not squeezed for a time long enough to remove the desired amount of liquid. When the draining of the sludge is assisted by squeezing the sludge between sandwiched endless conveyor belts, the smaller sludge particles tend to pass through the conveyor belts. Also, if the sandwiched conveyor belts are run through a sinuous path so as to create shear forces in the sludge filter cake, the shifting of the filter cake tends to break the filter cake and allow some of the finer solids to pass with the liquid through the filter belts. On the other hand, after a high percentage of the liquid has been removed from the sludge cake, the sludge cake can be squeezed and higher compression forces can be applied to the sludge cake to remove more liquid substantially without hazard of allowing the small solid particles to pass through the filter belt, since the reduced quantity of water present in the sludge cake is insufficient to wash out the small particles.\nOnce the sludge has been dewatered, it is desirable to remove the sludge from the pores of the filter belts, so that a porous belt is available on the next cycle of the belt for receiving the sludge slurry. The cleaning of the filter belts is usually accomplished by directing a continuous flow of water against the belt, and requires the excessive use of water.\nDewatered sludge can be used for various purposes, such as for a landfill as a fertilizer, or as a fuel. If the dewatered sludge has a low liquid content, it weighs less for a given volume, which means that it can be transported at a lower cost, and the dry product can be better used in subsequent processes since the energy required to extract the residual water would be less and the dryer the sludge, the better the sludge sustains combustion. Also, a dryer sludge makes a better landfill."}
{"text": "In communications networks, there may be a challenge to obtain good performance and capacity for a given communications protocol, its parameters and the physical environment in which the communications network is deployed.\nFor example, one parameter in providing good performance and capacity for a given communications protocol in a communications network is the availability of accurate system time information.\nFor example, accurate system time information may be required in radio-based communications networks in order to facilitate good sleep opportunities for both wireless devices and network nodes, saving significant power on both sides, without introducing any ambiguities in the communication links between the wireless devices and the network nodes. In general terms, time information is to the wireless device provided by the network node in the serving cell of the wireless device. By maintain per cell time information the need to synchronize cells is removed, simplifying deployment of radio-based communications networks.\nTime information is provided in system broadcasts and the frequency and granularity of the information is a trade-off between sleep efficiency of the wireless devices, the maximum sleep time of the wireless devices, and system overhead. For example, in telecommunications networks based on the long term evolution (LTE) standards the system frame number is provided in the master information block (MIB) which is transmitted once every 10 ms. The time granularity is 10.24 seconds, or 1024 frames, limited by the number of bits used in the MIB to signal the frame number.\nIt has been considered to extend the time information in communications networks to enable longer sleep-times for the wireless device, for example by reserving additional bits for signaling or to introduce new messages for providing a secondary timing, differentiating between cycles of the existing timing.\nHowever, in some communications networks, such as in some multi-layer multi technology environments, the ability of network energy saving by network node dormancy and overhead saving is reduced by the need to provide frequent time-information to enable efficient sleep of the wireless devices. The network node still frequently needs to provide time information in order for the wireless device to quickly (re-)acquire timing information.\nHence, there is still a need for an improved provision of time information to a wireless device."}
{"text": "In recent years, along with the popularization of personal computers, it becomes easy to input an image using an image input device such as a digital camera, color scanner, or the like, to display and confirm that image using an image display device such as a CRT, LCD, or the like, and to output the image via an image output device such as a color printer or the like. In this case, in order to correct any color appearance difference due to the difference between color gamuts of the display and output devices, a color matching process (gamut mapping) is normally executed. That is, the color appearance difference between the devices is absorbed by the color matching process.\nAs one of color matching processes, a method using an ICC profile is known. Note that the ICC profile is a profile based on a format defined by International Color Consortium.\nAs shown in FIG. 1, the process using an ICC profile converts the RGB values of an image obtained by the input device into CIEXYZ or CIEL*a*b* values as a device-independent color space using a source profile, and converts the CIEXYZ or CIEL*a*b* values into the device RGB values of the output (display) device using a destination profile. The source profile describes an RGB to XYZ conversion matrix or an RGB to Lab conversion lookup table (LUT) that represents the color reproduction character of the input device. On the other hand, the destination profile describes a Lab to device RGB LUT according to a specific format. The LUT of this destination profile has the following format. That is, grid point data obtained by evenly n-slicing (n is an integer; to divide into n−1 slices) an L* value region 0≦L*≦100, an a* value region −128≦a*≦127, and an b* value region −128≦b*≦127 are defined, and the LUT describes device RGB values obtained by mapping these L*a*b* values onto the color gamut of the output device.\nIn this manner, in the processing system of the ICC profile, color matching information, i.e., gamut mapping information is described in the LUT in the destination profile together with color gamut information of the output device.\nOn the other hand, a color appearance model that models the human visual character with higher precision than the conventional CIELAB color space has been announced. The color values (e.g., J*a*b* values) of this color appearance space can be calculated from CIEXYZ values. By making gamut mapping using these J*a*b* values, color matching more suited to the human visual character can be done.\nWhen an estimation of color matching using the color appearance space realized by the processing system of the ICC profile is made, a new destination profile must be generated from a destination profile of the ICC profile. Note that this generation method is the invention to explain later, however, this generation method is explained in here to clear a subject of the invention.\nThe destination profile is embedded with the LUT in the predetermined format that describes the relationship between the L*a*b* values and RGB values (DevRGB). In other words, in order to convert color matching information using the color appearance space into a destination profile, the following processes are required: (i) the L*a*b* grid points in the predetermined format are converted into J*a*b* values using color appearance conversion formulas, (ii) the J*a*b* values are mapped into the color gamut of the output device, (iii) the mapped J*a*b* values are converted into the device RGB, and (iv) the device RGB values are stored in an LUT. The flow of calculations is L*a*b*→J*a*b*→J*′a*′b*′→DevRGB, as shown in FIG. 2.\nHowever, the conversion formula required to obtain the color values J*a*b* of the color appearance space does not support colors outside of the visible region. In the color appearance conversion, for example, color adaptation conversion, cone response conversion, or opponent-colors response conversion is made so the color values approximate the human visible character. The conversion method of the color adaptation conversion is linear conversion represented by matrix conversion given by:\n                              A          ′                =                              (                                                                                a                    00                                                                                        a                    01                                                                                        a                    02                                                                                                                    a                    10                                                                                        a                    11                                                                                        a                    12                                                                                                                    a                    20                                                                                        a                    21                                                                                        a                    22                                                                        )                    ⁢          A                                    (        1        )            The cone response conversion also performs, in addition to the matrix conversion, an exponentiation given by:\n                              A          ′                =                                                            b                ⁡                                  (                                      cA                    /                    d                                    )                                            e                                      f              +                                                (                                      cA                    /                    d                                    )                                e                                              +          g                                    (        2        )            where a00 to a22, and b to g are constants (for 0≦e≦1), A is a value before conversion, and A′ is a value after conversion.\nThe opponent-colors response conversion converts the values after cone response conversion into coordinate J*a*b* values suited to the human visual character, and linear conversion and triangular function conversion are used as its representative conversion formulas.\nAccording to these conversion formulas, when XYZ values indicating a color outside of the visible region are input, the color adaptation conversion (equation (1)) as the linear conversion can be calculated, but calculation errors may occur when the exponentiation like equation (2) is calculated. This is for the following reason. That is, the constants in equations (1) and (2) are set so that (cA/d) in equation (2) assumes a positive value in case of the color within the visible region. However, (cA/d) often assumes a negative value when a color outside of the visible region is input, thus disabling calculations.\nMore specifically, as shown in FIG. 3, of the L*a*b* grit points defined on the CIELAB space, a color (e.g., L*=0, a*=−128, b*=−128) outside of the visible region cannot be converted into J*a*b* values, and device RGB values corresponding to that grid point cannot be calculated. In practice, when color matching is done by the processing system of the ICC profile, it is impossible that an input image includes a color outside of the visible region, but a color in proximity of a boundary of the visible region often exists. In order to interpolate such color using an LUT, since not only grid points within the visible region but also those outside of the visible region are used, appropriate device RGB values must be set for the grid points outside of the visible region.\nAs a color matching method using the ICC profile and color appearance space, a technique disclosed in Japanese Patent Laid-Open No. 2000-40140 (corresponding to U.S. Pat. No. 7,027,067, issued Apr. 11, 2006) has been proposed. However, the method disclosed in this reference is different from the method described in FIG. 2. The technique disclosed in this reference loads the color gamuts of respective devices from the source profile and destination profile described in the ICC profile format into a dedicated apparatus, and attains gamut mapping by converting the respective color gamuts into J*a*b* values using a CIECAM97 conversion formula as one of color appearance conversion formulas. When this technique is used, a new destination profile that records gamut mapping information on the J*a*b* space need not be generated, and the same profiles as the conventional ones need only be used. Hence, no problems caused by impossibility of color conversion as described above occur. However, since this technique requires a special apparatus upon color matching, it cannot be used in a general image manipulation application such as PhotoShop®. Since color matching is done using the dedicated apparatus, gamut mapping information described in the ICC profile is completely ignored.\nIn the processing system of the ICC profile, it is required to adopt gamut mapping using the color appearance space. Not only when an ICC profile is generated but also when a process is made on the color appearance space, it is inconvenient that colors outside of the visible region are not supported."}
{"text": "The present invention relates generally to freight shipping and, more particularly, to a method and apparatus for tracking air cargo.\nFreight or cargo is shipped using any of a variety of methods, the method or methods used for any one shipment being selected on the basis of cost, delivery schedule, and pick-up and drop-off locations. Furthermore, if the goods are perishable or fragile, additional shipping precautions may be required.\nDuring shipment, it is extremely difficult and typically very labor intensive to determine the exact location of the goods being shipped. Such tracking often requires multiple communications via phone, fax, e-mail, or other means between the shipper, forwarder, and the various carrier personnel. Location determination may even require physically searching a cargo storing warehouse, trucking warehouse, airport cargo terminal, or other location.\nShipment tracking is complicated by a number of factors. First, a variety of parties typically handle a piece of cargo during shipping, the parties potentially working for multiple independent companies. For example, a single shipment may require handling by multiple trucks and/or trucking companies and utilize multiple airline flights and/or airline companies. Additionally, a variety of cargo handlers/loaders are required to move the cargo between trucks/flights/carriers, staging areas, and temporary storage locations. Second, there are often last minute changes in shipping, for example due to a piece of freight being off-loaded from a scheduled flight in order to make room for a more valuable or time sensitive piece of freight. Such changes are typically made at the discretion of the carrier without notifying either the shipper or the forwarder. In this situation the shipper or forwarder may not be aware of the change until the freight does not arrive on the scheduled flight or at the expected time. As a result, the shipper is likely to expend additional time attempting to locate the delayed freight.\nA number of carriers have attempted to provide improved tracking information by allowing at least some of the larger forwarders to directly view their loading manifests, for example via a Web site. Although this approach can provide some tracking information, it requires carriers to accurately update the status of cargo at each step of cargo movement. As a consequence, this approach is labor intensive and thus costly. Additionally, the cargo status may not be updated at each step, for example, the carrier may only update the status when the cargo actually is loaded or unloaded from a plane, not when the cargo is simply moved from a storage facility to a loading site. Lastly, as cargo status information relies on human operators inputting the data, the data may be incorrect or delayed.\nSome shipping integrators have begun to use bar code tracking as a means of providing cargo tracking. Typically each piece of cargo is immediately tagged with a bar code upon receipt by the integrator. Then at each step of shipping, the handler scans the bar code, thereby immediately logging the location of the cargo as well as the time that the bar code was scanned into a manifest readable by the shipper via a Web site or other means. Although this approach reduces the risk of human error through the use of bar codes, the approach is still labor intensive as human operators are required to input the data. Additionally, the provided information is still only as accurate as the last bar code scan.\nAccordingly, what is needed in the art is a system that autonomously tracks cargo during shipping, providing the shipper and/or forwarder with accurate, timely, cargo status. The present invention provides such a system.\nAn autonomous cargo tracking system and method of using same are provided. The system is comprised of a communication and sensor unit that is affixed to the container to be tracked, a ground system that coordinates communications between the unit and the users, and one or more satellite systems that provide communication capabilities as well as position information. The system can operate either as a unidirectional system or as a bi-directional system. As a unidirectional system data is sent either directly to the users or, more preferably, sent indirectly to the users via the ground system. As a bi-directional system users are able to both request and receive data from the unit, either directly or indirectly through the ground system. Preferably users communicate with the ground system using an Internet Web site.\nAccording to the invention, a position sensing and communication (PSC) unit is affixed to or integrated into a shipping container, shipping pallet, cargo net, or cargo unit load device (ULD). The PSC unit provides a means of communicating with a ground system and determining location coordinates. Additionally, the PSC unit can include one or more sensor modules that can be used to provide additional information such as environmental conditions, container status, rate and/or direction of movement, etc. Once activated, the identity of a specific PSC unit is linked to or associated with the specific cargo to which it is coupled.\nThe PSC unit, using a GPS satellite constellation or other means, obtains position coordinates in response to a triggering event, the type of triggering event depending upon the configuration of the PSC unit and/or the system. Triggering events include detected movement, detected movement meeting preset criteria, a change in a predefined environmental condition (e.g., humidity, vibrational level, pressure, temperature, etc.), lapsing of a preset time period, and/or when polled by the ground system or a user.\nIn the preferred embodiment position coordinates, once determined, are sent to a ground system for translation into user recognizable locations, the translated coordinates provided either textually (i.e., as a location description) or graphically (i.e., as a map with a location indicator). The ground system sends the location information to a user, either in response to a request for the information by the user or automatically according to a predefined set of rules.\nPreferably the communication portions of the PSC unit are placed in standby mode when the container and the associated PSC unit are loaded into the cargo bay of an aircraft, thus insuring that the unit does not interfere with the aircraft\"\"s avionics systems. In one embodiment, container loading is determined on the basis of the container being raised more than five feet, as determined by an on-board accelerator module, and at a rate commensurate with a cargo loader. Preferably container loading is verified, for example by comparing the container\"\"s coordinates with an airport data base. In an alternate embodiment, container loading is determined by monitoring the electromagnetic spectrum of its surroundings, specifically monitoring for the radiated emissions of the aircraft power supply which are most prevalent at 400 Hertz. In another alternate embodiment, the PSC unit includes a RF detector designed to detect non-aircraft originating RF signals, the receipt of which is interrupted when the container and PSC unit are loaded into the aircraft. In yet another alternate embodiment, the PSC unit\"\"s sensor module emits one or more pulses of radiation, preferably ultrasonic radiation, and monitors the return of the radiation. The timing of the return can be used to detect the proximity of adjacent surfaces, such as the walls of the aircraft.\nA further understanding of the nature and advantages of the present invention may be realized by reference to the remaining portions of the specification and the drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a push handle switch lampholder, and more particularly to the push handle switch lampholder that pushes a lower conductive plate to move like a see-saw by pushing a push handle to move back and forth and determines whether or not to pass a current to an upper conductive plate, so as to control the ON or Off of a light emitting element. In the meantime, a set of hooks is provided for connecting a power cable for an electric connection without the need of tying an UL knot.\n2. Description of the Related Art\nWith reference to FIG. 1 for a conventional push handle switch lampholder 1, the push handle switch lampholder 1 comprises a cylindrical shell 11, an insulating plate 12, an insulating base 13 and a push handle switch 14, wherein the cylindrical shell 11 is electrically conductive and provided for screwing and connecting a light emitting element such as a light bulb (not shown in the figure), and a first electrode of the light emitting element is contacted with cylindrical shell 11 to electrically couple to a negative electrode of the power supply such as a utility power supply (not shown in the figure). The insulating plate 12 is installed at the bottom of the cylindrical shell 11, and a connecting groove 120 is formed thereon. The insulating base 13 is installed under the insulating plate 12 and includes a containing groove 130 concavely formed on the insulating base 13 for embedding a negative electrode plate 131 and a positive electrode plate 132, and an end of the negative electrode plate 131 is extended to the top of the insulating base 13 and abutted against the bottom of the insulating plate 12, and ends of the cylindrical shell 11, the insulating plate 12, the negative electrode plate 131 and the insulating base 13 are pivotally coupled, so that the negative electrode plate 131 and the cylindrical shell 11 are electrically connected. The positive electrode plate 132 includes a lower conductive plate 132a, at least two conductive plates 132b and an upper conductive plate 132c, and the lower conductive plate 132a is embedded into the insulating base 131. The conductive plates 132b are contained in both sides of the containing groove 130 respectively, and one of the conductive plates 132b proximate to the lower conductive plate 132a is electrically coupled to the lower conductive plate 132a, and the other of the conductive plates 132b is electrically coupled to the upper conductive plate 132c, and the upper conductive plate 132c is passed out of the connecting groove 120 and contacted with the second electrode of the light emitting element. A set of plug slots 133 is formed on an external surface of the insulating base 13 and their interior is corresponding to the negative electrode plate 131 and lower conductive plate 132a separately.\nThe switch 14 includes a push handle 141, a spring 142 and a metal plate 143, and the push handle 141 is contained in containing groove 130 for performing a pushing action reciprocally in the containing groove 130, and a middle section of the push handle 141 is hooked to an end of the spring 142. The metal plate 143 is substantially C-shaped, and an internal edge of the lower section is latched to the middle section of the push handle 141, and an internal edge of the upper section is hooked to the other end of the spring 142, so that when the push handle is pushed, the push handle is driven by an elastic force of the spring 142 to tilt towards one of the ends of the push handle 141, and the metal plate 143 is contacted with the conductive plates 132b on both sides respectively (as shown in FIG. 2). When the push handle 141 drives the spring 142 to link and move the metal plate 143 to tilt towards another end of the push handle 141, the metal plate 143 is released from the contact of the conductive plates on both sides.\nIn FIG. 3, both bare ends of the power cable 15 are plugged into two plug slots 133 respectively, so that the two bare ends of the power cable 15 push the negative electrode plate 131 and the lower conductive plate 132a to connect with each other for an electric connection. Therefore, the push handle 141 can push the metal plate 143 to contact the conductive plate 132b on both sides the same time to achieve an ON state, so as to electrically conduct the positive terminal of the power source to light up the light emitting element. On the other hand, the push handle 141 can push the metal plate 143 to disconnect the conductive plate 132b on both sides to achieve an OFF state, so as to turn off the light emitting element.\nThe connection between the push handle switch lampholder 1 and the power cable 15 relies on the negative electrode plate 131 and the lower conductive plate 132a to connect both bare ends of the power cable 15 only, so that when the power cable 15 is pulled accidentally or intentionally, the bare ends of the power cable 15 may be separated from the negative electrode plate 131 and/or the lower conductive plate 132a to result in a circuit disconnection. In the use of the conventional push handle switch lampholder 1, an UL knot 151 is tied at a position near the bare end of the power cable 15, and then both bare ends are plugged into two plug slots 133 respectively to achieve the effect of connecting the negative electrode plate 131 and the lower conductive plate 132a for an electric connection, and the UL knot 151 is intended for increasing the tensile resistance between the switch lampholder 1 and the power cable 15.\nThe conventional push handle switch lampholder 1 has the following drawbacks:\n(1) Difficult Assembly: Since the lampholder 1 uses a number of conductive plates for conducting current to the light emitting element, therefore the overall assembling operation is complicated and tedious, and the position may be shift easily by an external force to affect the effect of the lampholder 1, so that the conventional lampholder 1 incurs a complicated assembling operation and the assembly may be loosened easily by external pushing or squeezing forces to affect the effect of use.\n(2) Large Volume: Since the whole switch 14 must be accommodated into the containing groove 130, therefore the insulating base 13 is limited and its volume cannot be reduced effectively. As a result, the cost for the molding process of the insulating base 13 is high, and the large volume will also affect the storage and transportation costs.\n(3) Unsmooth Pushing Operation: When a user pushes the push handle 141, a force greater than the elastic force of the spring must be applied to the push handle 141 before the metal plate 143 can be driven to contact with or separate from the positive electrode plates 132 on both sides. If the elastic force of the spring 142 is too large, then the user has to push the push handle 141 inconveniently. If the elastic force of the spring 142 is too small, then the push handle 141 will be too loose to cause a wrong movement. In addition, the spring 142 may become elastically fatigue after a long time of use, so as to affect the operation and effect of the lampholder 1.\n(4) Short Service Life: Since the metal plate 143 and the conductive plates 132b on both sides are contacted to allow a current flow, and the contact position is a point or a line, the metal plate 143 produces a high temperature and becomes deteriorated or broken after a long time of use, so as to shorten the service life of the lampholder 1.\n(5) Complicated Tie of UL Knot: The process of tying an UL knot 151 at a position proximate to a bare end of the power cable 15 is troublesome, not only increasing the manufacturing time and cost, but also increasing the storage and transportation costs since the UL knot 151 has a relatively large volume and it is difficult to reduce the total volume of the switch lampholder 1. Particularly, the UL knot increases the tensile resistance between the switch lampholder 1 and the power cable 15, but fails to assure a secured connection between the switch lampholder 1 and the power cable 15 when a strong force is applied."}
{"text": "Compounds which have retinoid-like activity are well known in the art and are described in numerous patents and scientific publications. It is generally accepted that pharmaceutical compositions having a retinoid-like compound are useful for treating and/or preventing skin-related diseases such as, but not limited to acne, actinic keratosis, psoriasis, eczema and atopic dermatitis. It is also known that they are useful to reverse or treat the effects of age and photo damage to the skin and to prevent and/or treat cancerous or precancerous conditions.\nOne of the most significant drawbacks associated with the use of retinoids, especially in the topical treatment of dermatological diseases, is local irritation. Retinoids or compounds having retinoid-like activity that combine good topical efficacy and cutaneous tolerability are not very common. Recently, the new drug adapalene (\"Differin\", CIRD Galderma) was reported to offer these favorable characteristics and it has been launched in several countries as a water-based gel formulation. ##STR2##\nU.S. Pat. No. 5,648,514 describes a series of substituted (5,6)-dihydronaphthalene derivatives having retinoid-like biological activity of the formula ##STR3##\nwherein\nR.sub.1 is hydrogen or alkyl of 1 to 10 carbons; PA1 R.sub.2 and R.sub.3 are hydrogen, or alkyl of 1 to 6 carbons and the substituted ethynyl group occupies either the 2 or the 3 position of the dihydronaphthalene nucleus; PA1 m is an integer having the value of 0-3; PA1 o is an integer having the value 0-3; PA1 Y is a phenyl group, or heteroaryl selected from a group consisting of pyridyl, thienyl, furyl, pyridazinyl, pyrimidinyl, pyrazinyl, thiazolyl, oxazolyl, and imidazolyl, said groups being optionally substituted with one or two R.sub.2 groups; PA1 A is (CH.sub.2).sub.n where n is 0-5, lower branched chain alkyl having 3-6 carbons, cycloalkyl having 3-6 carbons, alkenyl having 2-6 carbons and 1 or 2 double bonds, alkynyl having 2-6 carbons and 1 or 2 triple bonds; PA1 B is hydrogen, COOH or a pharmaceutically acceptable salt thereof, COOR.sub.8, CONR.sub.9 R.sub.10, CH.sub.2 OH, CH.sub.2 OR.sub.11, CH.sub.2 OCOR.sub.11, CHO, CH(OR.sub.12).sub.2, CHOR.sub.13 O, COR.sub.7, CR.sub.7 (OR.sub.12).sub.2, CR.sub.7 OR.sub.13 O, or trilower alkylsilyl, where R.sub.7 is an alkyl, cycloalkyl or alkenyl group containing 1 to 5 carbons, R.sub.8 is an alkyl group of 1 to 10 carbons or trimethylsilylalkyl where the alkyl group has 1 to 10 carbons, or a cycloalkyl group of 5 to 10 carbons, or R.sub.8 is phenyl or lower alkylphenyl, R.sub.9 and R.sub.10 independently are hydrogen, an alkyl group of 1 to 10 carbons, or a cycloalkyl group of 5-10 carbons, or phenyl or lower alkylphenyl, R.sub.11 is lower alkyl, phenyl or lower alkylphenyl, R.sub.12 is lower alkyl, and R.sub.13 is divalent alkyl radical of 2-5 carbons; and PA1 R.sub.22 is hydrogen, alkyl of 1 to 10 carbons, fluoro-substituted alkyl of 1 to 10 carbons, alkenyl of 2 to 10 carbons and having 1 to 3 double bonds, alkynyl having 2 to 10 carbons and 1 to 3 triple bonds, carbocyclic aryl selected from the group consisting of phenyl, C.sub.1 -C.sub.10 -alkylphenyl, naphthyl, C.sub.1 -C.sub.10 -alkylnaphthyl, phenyl, C.sub.1 -C.sub.10 alkyl, naphthyl-C.sub.1 -C.sub.10 alkyl, C.sub.1 -C.sub.10 -alkenylphenyl having 1 to 3 double bonds, C.sub.1 -C.sub.10 -alkynylphenyl having 1 to 3 triple bonds, phenyl-C.sub.1 -C.sub.10 alkenyl having 1 to 3 double bonds, phenyl-C.sub.1 -C.sub.10 alkynyl having 1 to 3 triple bonds, hydroxyalkyl of 1 to 10 carbons, hydroxyalkynyl having 2 to 10 carbons and 1 to 3 triple bonds, acyloxyalkyl of 1 to 10 carbons or acyloxyalkynyl of 2 to 10 carbons and 1 to 3 triple bonds, where the acyl group is represented by COR.sub.14, CN, CON(R.sub.1).sub.2, (CH.sub.2).sub.p CO.sub.2 R.sub.8 where p is an integer between 0 to 10, or R.sub.22 is aminoalkyl or thioalkyl of 1 to 10carbons, or a 5 or 6 membered heteroaryl group optionally substituted with a C.sub.1 to C.sub.10 alkyl group and having 1 to 3 heteroatoms, said heteroatoms being selected from a group consisting of O, S, and N, or R.sub.22 is represented by (CH.sub.2).sub.p XR.sub.1 or by (CH.sub.2).sub.p NR.sub.1 R.sub.2 ; where X is O or S, the R.sub.14 group is hydrogen, alkyl of 1 to 10 carbons, alkenyl of 2 to 10 carbons and having 1 to 3 double bonds, alkynyl having 2 to 10 carbons and 1 to 3 triple bonds, carbocyclic aryl selected from the group consisting of phenyl, C.sub.1 -C.sub.10 -alkylphenyl, naphthyl, C.sub.1 -C.sub.10 -alkylnaphthyl, phenyl-C.sub.1 -C.sub.10 alkyl, or naphthyl-C.sub.1 -C.sub.10 alkyl. PA1 X.sub.1 is [C(R.sub.1).sub.2 ].sub.n where R.sub.1 is independently H or alkyl of 1 to 6 carbons, and n is an integer between 0 and 2; ##STR5## PA1 R.sub.2 is hydrogen, lower alkyl of 1 to 6 carbons, F, Cl, Br, I, CF.sub.3, fluoro substituted alkyl of 1 to 6 carbons, OH, SH, alkoxy of 1 to 6 carbons, or alkylthio of 1 to 6 carbons; PA1 R.sub.3 is hydrogen, lower alkyl of 1 to 6 carbons or F; PA1 m is an integer having the value of 0 to 3; PA1 o is an integer having the value of 0 to 3; PA1 Y is a phenyl or naphthyl group, or heteroaryl selected from a group consisting of pyridyl, thienyl, furyl pyridazinyl, pyrimidinyl, pyrazinyl, thiazolyl, oxazolyl, imidazolyl and pyrrazolyl, said phenyl and heteroaryl groups being optionally substituted with one or two R.sub.2 groups, or PA1 when Z is --(CR.sub.1.dbd.CR.sub.1).sub.n' -- and n' is 3, 4 or 5 then Y represents a direct valence bond between said (CR.sub.2.dbd.CR.sub.2).sub.n' group and B; PA1 A is (CH.sub.2).sub.q where q is 0-5, lower branched chain alkyl having 3-6 carbons, cycloalkyl having 3-6 carbons, alkenyl having 2-6 carbons and 1 or 2 double bonds, alkynyl having 2-6 carbons and 1 or 2 triple bonds; PA1 B is hydrogen, COOH or a pharmaceutically acceptable salt thereof, COOR.sub.8, CONR.sub.9 R.sub.10, CH.sub.2 OH, CH.sub.2 OR.sub.11, CH.sub.2 OCOR.sub.11, CHO, CH(OR.sub.12).sub.2, CHOR.sub.13 O, --COR.sub.7, CR.sub.7 (OR.sub.12).sub.2, CR.sub.7 OR.sub.13 O, or Si(C.sub.1-6 alkyl).sub.3, where R.sub.7 is an alkyl, cycloalkyl or alkenyl group containing 1 to 5 carbons, R.sub.8 is an alkyl group of 1 to 10 carbons or (trimethylsilyl)alkyl where the alkyl group has 1 to 10 carbons, or a cycloalkyl group of 5 to 10 carbons, or R.sub.8 is phenyl or lower alkylphenyl, R.sub.9 and R.sub.10 independently are hydrogen, an alkyl group of 1 to 10 carbons, or a cycloalkyl group of 5-10 carbons, or phenyl or lower alkylphenyl, R.sub.11 is lower alkyl, phenyl or lower alkylphenyl, R.sub.12 is lower alkyl, and R.sub.13 is divalent alkyl radical of 2-5 carbons; and PA1 R.sub.14 is (R.sub.15).sub.r -substituted alkyl of 1-6 carbons, (R.sub.15).sub.r -substituted alkenyl of 1-6 carbons and 1 or 2 double bonds, (R.sub.15).sub.r -substituted alkynyl of 1-6 carbons and 1 or 2 triple bonds, (R.sub.15).sub.r -phenyl, (R.sub.15).sub.r -naphthyl, or (R.sub.15).sub.r -heteroaryl where the heteroaryl group has 1 to 3 heteroatoms selected from the group consisting of O, S and N, r is an integer having the values of 0-5, and PA1 R.sub.15 is independently H, F, Cl, Br, I, NO.sub.2, N(R.sub.8).sub.2, N(R.sub.8)COR.sub.8, NR.sub.8 CON(R.sub.8).sub.2, OH, OCOR.sub.8, OR.sub.8, CN, COOH, COOR.sub.8 and alkyl group having 1 to 10 carbons, fluoro substituted alkyl group having 1 to 10 carbons, an alkenyl group having 1 to 10 carbons and 1 to 3 double bonds, alkynyl group having 1 to 10 carbons and 1 to 3 triple bonds, or a (trialkyl)silyl or (trialkyl)silyloxy group where the alkyl groups independently have 1 to 6 carbons.\nThat disclosure is specifically limited to the ethynyl linker. The two compounds disclosed in the present invention have the ethenyl linker. The substituent in position 8 (R.sup.22) is defined as being, among others, an alkenyl group of 2 to 10 carbons and having 1 to 3 double bonds or an alkynyl group having 2 to 10 carbons and 1 to 3 triple bonds. This general definition does not specify the direct attachment to the dihydronaphthalene nucleus at position 8.\nPublished PCT patent application WO 97/48672 discloses a series of 5,6-dihydronaphthalene derivatives having retinoid and/or retinoid antagonist-like activity of the formula ##STR4##\nwherein\nTwo of the agents included within the scope of WO 97/48672 are the compounds having the structural formula I and II ##STR6##\nThese compounds show very potent activity against dermatological diseases. Unexpectedly, the present inventors have discovered that these compounds show also a substantially reduced irritancy profile. While WO 97/48672 generally discloses these compounds as having retinoid-like activity as agents for treating skin-related diseases, there is no specific disclosure of compounds I and II or their unexpectedly low-irritation profile, which are the subject of the present invention.\nIn another aspect, the present inventors have also found that compounds I and II are effective tumor inhibiting agents, and thus are useful in human and/or veterinary medicine. WO 97/48672 generally discloses these compounds having retinoid-like activity and as agents for preventing or treating cancerous and precancerous conditions. The assay supporting this statement is a measure of the inhibition of 12-O-tetradecanoylphorbol-13-acetate (TPA) induction of ornithine decarboxylase (ODC) in mouse epidermis by certain compounds disclosed in application. TPA-induced ODC activity is known to occur at pre-malignant stages. Although the correlation of this assay with cellular anti-proliferation is well established, there is no evidence in the WO 97/48672 patent application of in vivo inhibition of fully established tumor growth by the disclosed compounds. The present inventors have found that the anti-proliferation activity of the compounds of this invention is translated to a potent inhibition of established tumor growth, equivalent to the most potent antitumor agents such as doxorubicin."}
{"text": "The present invention relates to a hermetic reciprocating compressor including a hermetic casing, a compressing portion having a cylinder and a piston reciprocating inside the cylinder and a motor for driving the piston.\nA hermetic compressor is generally employed in a cooling system such as a refrigerator or an air conditioner, so as to compress a gaseous refrigerant received from an evaporator and supply the compressed refrigerant to a condenser.\nFIGS. 3 and 4 show front and side sectional views of a conventional hermetic reciprocating compressor, respectively. As shown in the drawings, the conventional compressor includes a casing 101 forming a closed internal space, a driving motor 110 installed inside the casing 101 and a compressing portion 120 which is driven by the driving motor 110 to compress a refrigerant. The driving motor 110 includes a stator 112, a rotor 111 rotatably installed inside the stator 112 and a crank shaft 117 fitted into the rotor 111 and rotating therewith while supported within a support member member. The compressing portion 120 includes a cylinder 113, a piston 123 reciprocating inside the cylinder 113 and a cylinder head 127. The piston 123 is connected to an eccentric portion 119 of the crank shaft 117 by a connecting rod 121 and reciprocates inside the cylinder 113 according to a rotational movement of the rotor 111, thereby inhaling and compressing the refrigerant. A suction muffler 141 is installed on the cylinder head 127 to guide the refrigerant to be compressed into an internal space 125 of the cylinder 113. A suction pipe 129 for transferring the refrigerant from an evaporator (not shown) to the compressor is connected to the suction muffler 141 after passing through a wall of the casing 101.\nReferring to FIG. 5, the suction muffler 141 has an internal space for receiving the refrigerant, an inlet 142 to which the suction pipe 129 is connected and an outlet which communicates with a refrigerant inlet 143 provided in the cylinder head 127. A coil spring 145 is installed between the inlet 142 of the suction muffler 141 and the suction pipe 129. One end of the coil spring 145 is fixedly inserted into the inlet 142 of the suction muffler 141 and the other end thereof is fitted outside the leading end of the suction pipe 129, so as to guide the refrigerant passing through the suction pipe 129 into the suction muffler 141.\nThe gaseous refrigerant from the evaporator flows into the suction muffler 141 via the suction pipe 129 and the coil spring 145, and is then supplied to the internal space 125 of the cylinder 113 through the cylinder head 127. On the other hand, the gaseous refrigerant supplied from the evaporator to the cylinder 113 contains liquid oil for lubrication and rust prevention for internal parts of the refrigerant circulation system. Since the refrigerant is vaporized in the evaporator by absorbing heat from the surroundings, whereas the oil maintains the liquid state due to its having a higher evaporation point than that of the refrigerant, the oil in the liquid state flows together with the gaseous refrigerant. The mixture of the liquid oil and the gaseous refrigerant contained in the internal space 125 of the cylinder 113 may damage the piston 123, the inner wall of the cylinder 113 or a valve plate (not shown) of the compressing portion 120, due to a liquid compression phenomenon of the liquid oil. Also, the liquid oil, having a relatively high specific volume, may obstruct the compression of the gaseous refrigerant, thereby decreasing the compression efficiency of the compressor."}
{"text": "Parametric sound is a fundamentally new class of audio, which relies on a non-linear mixing of an audio signal with an ultrasonic carrier. One of the key enablers for this technology is a high-amplitude, efficient ultrasonic source, which is referred to here as an emitter or transducer. Ultrasonic emitters can be created through a variety of different fundamental mechanisms, such as piezoelectric, electrostatic, and thermoacoustic, to name a few. Electrostatic emitters are generally capacitive devices consisting of two conductive faces with an air gap, where at least one of the conductive faces has a texture that is critical to the functionality of the emitter.\nNon-linear transduction results from the introduction of sufficiently intense, audio-modulated ultrasonic signals into an air column. Self-demodulation, or down-conversion, occurs along the air column resulting in the production of an audible acoustic signal. This process occurs because of the known physical principle that when two sound waves with different frequencies are radiated simultaneously in the same medium, a modulated waveform including the sum and difference of the two frequencies is produced by the non-linear (parametric) interaction of the two sound waves. When the two original sound waves are ultrasonic waves and the difference between them is selected to be an audio frequency, an audible sound can be generated by the parametric interaction.\nParametric audio reproduction systems produce sound through the heterodyning of two acoustic signals in a non-linear process that occurs in a medium such as air. The acoustic signals are typically in the ultrasound frequency range. The non-linearity of the medium results in acoustic signals produced by the medium that are the sum and difference of the acoustic signals. Thus, two ultrasound signals that are separated in frequency can result in a difference tone that is within the 60 Hz to 20,000 Hz range of human hearing.\nLike conventional audio systems, ultrasonic audio systems can suffer from distortion caused by the phenomenon known as clipping. Clipping is a form of waveform distortion that cuts off the peaks and troughs of the audio waveform when the signal is driven beyond the capacity of the amplifier. Clipping occurs when an audio signal exceeds a maximum value allowed by the audio system. The signal beyond the capability of the system is cut off, or clipped, resulting in distortion of the audio signal and the creation of nonlinear distortion products, such as unwanted harmonics related to the input.\nDigital signal processing (DSP) has been applied extensively to audio, including to perform functions such as compression, equalization, surround sound, pitch control, etc. Many of these algorithms can apply gain to an input signal and the resulting output can exceed the maximum signal level allowed by digital to analog conversion after processing. The resulting signal is ‘clipped’ whereby any signal attempting to exceed the maximum value is reduced, or clipped at the maximum value. Similarly, there is a minimum value whereby signals below this value are raised. To illustrate this, consider a system having its output limited to a range between −1 to 1 volts. Any value below −1 is raised to −1, and any value above 1 is dropped to 1. An example of this clipping is shown in FIG. 1. As seen in this example, an 1 kHZ input tone with an amplitude of 2 volts peak-to-peak (top of Figure) is clipped at +/−1 volt (bottom of Figure).\nThe resulting spectrum is shown in FIG. 2. Before clipping the spectrum is a pure tone with one frequency component (top of Figure). After clipping, odd (3rd, 5th, 7th, etc.) appear and significantly color the audio experienced (bottom of Figure)."}
{"text": "The present invention relates to a composition of vanadyl sulfate, alpha-lipoic acid, taurine and chromium carnitine to augment or normalize cellular tyrosine phosphorylation responses to insulin, thereby normalizing insulin activity, decreasing catabolic activity in the diabetic state, decreasing serum lipoproteins and maintaining normal body weight in healthy humans and animals.\nThe actions of insulin are initiated by its binding to insulin receptor, a disulfide-bonded heterotetrameric membrane protein (FIG. 1). Insulin binding causes conformational changes in insulin receptor that lead to autophosphorylation of the receptor and activation of the receptor\"\"s intrinsic tyrosine kinase activity. The function of the tyrosine kinase of insulin receptor is essential for the biological effects of insulin. After autophosphorylation, the insulin receptor can phosphorylate the tyrosine residues of several substrates, including the insulin receptor substrate (IRS) proteins, which in turn can activate downstream signaling molecules in hepatic, muscle and fat cells.\nA protein tyrosine phosphatase termed PTP-1B can terminate the signaling cascade initiated by insulin by removing the phosphate from the insulin receptor. Studies indicate that increased expression of PTP-1B in mice gives rise to a form of insulin resistance termed type 2 non-insulin dependent diabetes mellitus (NIDDM). NIDDM is a complex disease characterized by progressive development of insulin resistance and defects in insulin secretion, which often leads to overt hyperglycemia. On the other hand, some data indicate that mice lacking the PTP-1 B gene are more sensitive to insulin\"\"s blood glucose-lowering effects than are control animals, and, these PTP-1B defective mice appear to be more resistant to weight gain when consuming a high-fat diet {Elchebly et al., 238 SCIENCE 1544-1548 (1999)].\nVanadate is known to be an inhibitor of many phosphatases. Some researchers have proposed that vanadate administration may mimic the activity of insulin. For example, U.S. Pat. No. 5,550,113 and U.S. Pat. No. 5,888,993 disclose vanadium salts for blood sugar regulation. U.S. Pat. No. 5,885,980 and U.S. Pat. No. 5,614,224 disclose vanadium salts for the treatment of diabetes.\nHowever, vanadyl sulfate is not widely used to mimic the activity of insulin because chronic administration of vanadate can be toxic. Vanadyl sulfate has been shown to induce chromosomal damage and mitotic recombination in the fruit fly. Vanadyl sulfate can induce pulmonary inflammation in rats when inhaled. In vitro studies of bovine papilloma virus DNA-tranfected C3H/10T1/2 cells have demonstrated that vanadyl sulfate possess tumor promotion activity.\nTo avoid the toxic build-up of vanadate, some researchers suggest administering it intermittently. For example, U.S. Pat. No. 5,730,988 to Womack (Mar. 24, 1998) discloses administration of nutritional supplements containing a source of vanadate in a xe2x80x9cPhase Ixe2x80x9d supplementation program and lipoic acid and Gymnema sylvesire in a xe2x80x9cPhase IIxe2x80x9d supplementation cycle.\nGiven the potential usefulness of vanadyl sulfate in the treatment of diabetes, it would be desirable to administer compositions containing vanadyl sulfate to diabetic patients, and to normal individuals wishing to minimize their risk of developing NIDDM, if the undesirable side effects of vanadate could be avoided.\nAlpha-lipoic acid has been studied as a potential treatment for improving glucose metabolism. U.S. Pat. No. 5,693,664 to Wessel et al. (Dec. 2, 1997) discloses the R-(+) enantiomer of alpha-lipoic acid and its metabolites for the treatment of diabetes. A formulation of DL-lipoic acid is described in U.S. Pat. No. 5,599,835 to Fischer (Feb. 4, 1997) for treatment of a metabolic aberration of the multi-enzyme complex of pyruvate dehydrogenase, which is symptomatic of diabetes mellitus.\nTaurine is a conditionally essential amino acid that is found in the tissues of most animal species. It is not incorporated into proteins, but is found free in many tissues. Taurine is involved in a number of physiological processes including bile acid conjugation, osmoregulation, detoxification of xenobiotics, cell membrane stabilization, modulation of cellular calcium flux, and modulation of neuronal excitability. Low levels of taurine have been associated with retinal degeneration, growth retardation, and cardiomyopathy. Taurine has been used clinically in the treatment of cardiovascular diseases, hypercholesterolemia, seizure disorders, ocular disorders, diabetes, Alzheimer\"\"s disease, hepatic disorders, cystic fibrosis, and alcoholism. Animal and human studies indicate that taurine supplementation is effective in alleviating some of the complications of insulin-dependant diabetes. Taurine has been found to influence blood glucose and insulin levels, as well as increasing glycogen synthesis, and it may also be involved in the functioning and integrity of pancreatic beta cells.\nTrivalent chromium is an essential trace element for normal carbohydrate metabolism and insulin sensitivity. Because of this biological activity, chromium supplementation has been studied as a potential therapy of insulin resistant states and dyslipidemias, and has been promoted as a health aid to the general population. Chromium has principally been studied for its possible benefits in improving blood sugar control in diabetic patients.\nChromium improves the glucose/insulin system in subjects with hypoglycemia, hyperglycemia, diabetes and hyperlipemia with no detectable effects on control subjects. Chromium improves insulin binding, insulin receptor number, insulin internalization, beta cell sensitivity and insulin receptor enzymes with overall increases in insulin sensitivity. There have been several studies involving chromium supplementation of subjects with NIDDM and/or lipemia and most have reported beneficial effects of chromium on the glucose/insulin system. In summary, chromium is involved in the control of the glucose/insulin system and the form of chromium is critical when evaluating the role of chromium in this system.\nSeveral forms of chromium have been studied with respect to glucose metabolism and impact on insulin resistance. Controversy exists as to which supplemental form of chromium is preferable; and, regarding insulin resistance, this controversy is likely to continue. The present invention provides a novel composition of chromium (+3) as chromium carnitine. This chelated form of chromium, administered in combination with vanadyl sulfate, lipoic acid, and taurine, functions coordinately to normalizes insulin signaling in liver, muscle and fat cells.\nHence, while certain compounds are thought to separately mimic some of the activities of insulin, a single non-toxic composition that provides a more complete range of insulin-like activities would be desirable. The ideal formulation would (i) increase cell sensitivity to insulin, (ii) maintain normal serum glucose, (iii) decrease hypertriglyceridemia and/or lipoproteins, and (iv) maintain normal body weight even with consumption of a high fat diet. These objects and other are achieved in the present invention.\nIn accordance with the present invention, an oral formulation is provided that comprises vanadyl sulfate, alpha-lipoic acid, taurine and chromium carnitine. The present supplement; (a) increases cell sensitivity to insulin, (b) maintains normal serum glucose, (c) decreases serum triglycerides and/or serum lipoproteins, and (d) maintains normal body weight even with consumption of a high fat diet. The supplement can be used by normal, healthy individuals as well as diabetics.\nThe present invention also provides a method of oral supplementation in animals that includes administering to an animal suffering symptoms of diabetes a composition including therapeutically effective amounts of vanadyl sulfate, alpha-lipoic acid, taurine and chromium carnitine and continuing administration of the composition until the symptoms are reduced.\nThe present invention further provides a method of oral supplementation in animals that includes administering to an animal a composition, which includes therapeutically effective amounts of vanadyl sulfate, alpha-lipoic acid, taurine and chromium carnitine, and continuing administration of the composition until normalization of serum glucose is achieved.\nThe present invention still further provides a method of oral supplementation in animals that includes administering to an animal a composition, which includes therapeutically effective amounts of vanadyl sulfate, alpha-lipoic acid, taurine and chromium carnitine, and continuing administration of the composition until normalization of serum triglycerides or lipoproteins is achieved.\nThe present invention additionally provides a method of treating obesity in an animal that includes administering to an animal a composition including therapeutically effective amounts of vanadyl sulfate, alpha-lipoic acid, taurine and chromium carnitine.\nAccording to the present invention, the animal is a mammal selected from the group consisting of humans, nonhuman primates, dogs, cats, horses or cattle.\nA pharmaceutically acceptable carrier may also be used in the present compositions and formulations. The formulations can be manufactured in any form known to the skilled artisan, for example as solid capsules, caplets, softgels, liquids, bars, or functional foods.\nOther features and advantages of the invention will become apparent from the following detailed description, taken in conjunction with the accompanying drawings, which illustrate, by way of example, the features of the invention."}
{"text": "This invention relates to a method and apparatus for adjusting the quantity of liquid, typically water, deposited on the surfaces of fine granular materials, and a method of preparing mortar or concrete by utilizing the treated granular material.\nFine aggregates comprising river or mountain sand or artificial fine particles are widely used to prepare cement mortar or lime stone type hydraulic mortar which is used to construct buildings or many other civil structures. When digging or crushing various ores utilized in metallurgy or ceramic industry and coal, fine particles or dust are formed. Further, depending upon the field of use, it is necessary to crush the ores or coal into granules having a predetermined size. When pulverizing or refining these substances or when using them for chemical reactions, fine granules in the form of sludge or the like are often formed. As is well known, these fine granules contain substantial quantity of water adhering to the surface thereof. This is true not only for river or mountain sand but also coal. Especially, in recent years, these substances are dug by using jet water so that the quantity of the deposited water is considerably large. Even with converter slag which is free from water when it is formed, water is used to collect it. Moreover, as these materials are stored in the outdoor they are wetted by rain, dew or snow. Such wet particles can not be used directly. For example, when sintering or converting these materials to coke, and even when they are directly charged to a furnace, it is necessary to preliminally dry them before actual use. This requires extra heat energy, i.e., fuel. As will be described later in more detail, when a fine aggregate composed of river or mountain sand is used for the preparation of mortar or cement, the quantity of the deposited or surface water is an important factor which influences the quality of the product. Although the composition and particle size of the sand also influences the quality of the product, so long as sand collected from the same source is used it is easy to utilize the sand having the same composition and particle size and it is rare to admix sands from different sources. When the sand contains particles of different size it is easy to classify them into fine, medium and coarse particles with a sieve and a small difference in the particle size does not result in a great difference in the quality of the product. However, the quantity of the surface water varies greatly depending upon the source, the methods of collecting, conveying and storing the sand. Moreover, the specific surface area of the fine particles of sand is large so that the relative quantity of the deposited water is large. Moreover, sand contains water in the interstices of the sand particles which varies from time to time depending upon weather conditions. More particulary, when sand from the same source is piled up on the ground, its water content varies at the top and the bottom, and in the morning and at noon. When preparing cement mortar or concrete by using a fine aggregate, the ratio of water to cement (W/C), the ratio of cement to fine aggregate (C/g) and the ratio of cement or sand to a coarse aggregate (S/G or C/G) have a great influence upon the strength of the resulting product, fluidity, moldability and workability. Thus, when an excessive quantity of water is incorporated, segregation and bleeding are inevitable, thus decreasing the mechanical strength of the product. On the other hand deficient water impairs moldability and pouring property, so that even when vibration or pressure is applied at the time of molding or pouring it is difficult to obtain a dense structure which also decreases the mechanical strength. As above described, notwithstanding the fact that it is essential to select adequate W/C, etc., as the quantity of the deposited water varies greatly and as it is difficult to simply and accurately measure the quantity of the surface water, it is difficult to realize ideal ratios W/C and S/C, etc. Although it has been proposed to completely dry the fine aggregate or to measure the weight thereof in water, such methods are not suitable at the field where a large quantity of sand is used. More particularly, the former method requires a large quantity of heat energy and time and the latter method requires a step of perfectly permeating water into sand to drive off air according to JIS (Japanese Industrial Standard) it is prescribed to immerse the sand in water for 24 hours and a step of draining the water contained in the sand."}
{"text": "The present invention generally relates to a screen saver apparatus and method, and in particular, to a screen saver apparatus and method based on monitoring the various changes of the signal representing an image being processed.\nScreen saver applications for computer systems are well known. The purpose of a screen saver application is to prevent the burning of a fixed pattern onto the phosphor screen of a display, when there is little or no screen activity. The common way of detecting this inactivity is to monitor whether there is any user keyboard input. In other words, a timer is started after each keyboard entry and if the timer expires before another entry, then the screen saver application will cause the screen to go blank or display a predetermined moving image, thereby preventing screen burn. The screen will be restored as soon as another keyboard entry is detected.\nOther prior screen saver systems take the approach of looking at the motion of the images being displayed. For example, Japanese Patent No. 6332418 A, assigned to Fujitsu General Ltd. of Japan, discloses a screen saver apparatus having a motion adaption circuit.\nThe motion adaption circuit examines a video signal and determines the motion characteristics of the corresponding image. If a still picture is detected, a timer is started. If the still picture is displayed after a predetermined time period, the power supply to the display is cut off, thereby preventing screen burn. These motion based systems are, however, complicated and typically requires costly digital processing circuitry to determine whether a motion is present. In addition, these systems enable the screen saver function based solely on the motion factor and nothing else, which may cause unnecessary blanking of the screen.\nAnother prior system is disclosed in Japanese patent publication number 09327031. This system determines the peak beam current value of one frame of video for a particular period. The system also detects the motion of an image based on an inter-frame of video reference. Thereafter, a brightness control unit reduces the beam current of the CRT gradually to a reference current value corresponding to the reference voltage value during the still picture display.\nThe above system, however, requires complicated and expensive circuitries. In particular, it detects the relative motion of an image by using a motion of an image by using a motion adaptive 3-dimension YC separation system. Theis requires costly components such as A/D converters and digital frame stores.\nThe present inventor recognizes the advantages of being able to implement a screen saver function for a multi-media display in a cost effective manner, preferably using the circuitry that is already being used for other functions.\nIn addition, the present inventor recognizes that there is no need to blank the screen, even if the screen is static for a period of time, if the image being displayed has characteristics that do not exceed certain thresholds. In other words, even if a screen image is static, there may be little chance of screen burn if the image being displayed is of low brightness and/or contrast, for example.\nTherefore, in one exemplary embodiment, a method and system for implementing a screen saver function is described. A peak value of a characteristic such as beam current magnitude of a video signal being processed is determined. An average value of the same characteristic of the video signal is also determined. Changes in this average value is additionally monitored. Therefore, the characteristic of the video signal is adjusted only if the peak value exceeds a first threshold and a change less than a second threshold is detected in the average value."}
{"text": "1. Field of the Invention\nThis invention relates to an audio reproducing apparatus and method, audio recording apparatus and method, audio recording and reproducing system, audio data transmission method, information receiving apparatus, and recording medium which are particularly suitable for use in a headphone stereo, and in particular to those configured to store in a portable headphone stereo desired audio data externally transmitted to the portable headphone stereo.\n2. Description of the Related Art\nAmong portable headphone stereo devices with which a user can enjoy reproduced music either during his walk or in an automobile, most widely distributed are those using an analog-recording magnetic tape such as compact cassette. A user of a portable headphone stereo using an analog-recording medium tape typically records in a compact cassette a desired piece of music selected from FM broadcasting programs, CDs (compact discs) or other music sources, and sets the compact cassette in his portable headphone stereo to enjoy audio reproduction either during his walk or in an automobile.\nHowever, as long as compact cassettes or other analog-recording magnetic tapes are used, great improvements in quality of sound are not prospective, and dubbing causes deterioration of signals. Moreover, it takes a long time for a user to record desired pieces of music from CD or other music sources. Additionally, with compact cassettes or other magnetic tapes, the access time is slow, and a user cannot quickly search out, reproduce or repeat his desired music.\nSome portable headphone stereos use CDs. Since CDs are media exclusive for recording, a user of portable headphone CD stereo buys CD recorded with his desired music, and places the CD in his portable headphone CD stereo to enjoy audio reproduction during his walk or in an automobile. CDs are digital recording media, the quality of sound is excellent. The access speed is high, and any desired music can be reproduced quickly. However, since CD headphone stereos are exclusively for reproduction, users cannot make their own music sources compiling their desired pieces of music. Moreover, CD headphone stereos are affected by vibrations, and sound is often skipped over due to external vibrations.\nAlso known are portable headphone stereo players using DAT (digital audio tape), NT (non-tracking tape(trade mark)) or other digital-recording magnetic tapes as their recording media. Deterioration of signals by dubbing hardly occurs in devices using DAT, NT or other digital magnetic tapes. DAT promises audio reproduction of a very high quality of sound. NT permits recording over a long time in a very small cassette. Here again, however, devices using magnetic tapes involve the same problems that the access speed is slow and it takes a long time for repeated reproduction or cue search of a desired piece of music.\nAnother type of portable headphone stereo players uses MD (multi-disk(trade mark)). MDs are media for both recording and reproduction, and users can record their desired pieces of music on MDs from CD or other music sources and can place them in their portable MD headphone stereos to enjoy reproduced music during a walk or in an automobile. MDs are digital-recording media, and promise an excellent quality of sound. The access speed is high, and any desired music can be reproduced quickly. As a shockproof memory is used, devices are less affected by external vibrations.\nAs reviewed above, various kinds of recording media have been proposed for use in portable headphone stereos. However, none of these recording media used in conventional portable headphone stereos meet all requirements from the viewpoints of easy use and quality of sound.\nDevices using analog-recording compact cassettes have a problem in quality of sound. Those using DAT, NT or other digital-recording magnetic tapes have a problem in access speed. CDs are exclusively for reproduction and weak against vibrations. MDs can be used for both recording and reproduction and can be miniaturized but, since the number of titles of MD on sale is not abundant, it is sometimes difficult for users to obtain their desired music. It takes time to dub CD or other music sources.\nAnother problem with MDs is the use of ATRAC (Adaptive Transform Acoustic Coding) compression and expansion ICs or other ICs developed for exclusive use in MDs and the use of a particular architecture, i.e., a particular circuit arrangement as a method for actually mounting the ICs in order to reduce the entire dimension and decrease the cost. Therefore, such devices cannot be used in different ways, namely, for recording and reproducing a recording medium other than MDs, for example.\nMoreover, tastes of users for music are diverged more and more, and the fashion of music largely changes in a short time. It is difficult to exactly catch the fashion of music with conventional music recording media.\nTaking it into account, the present Applicant previously proposed a portable headphone stereo disclosed in Japanese Patent Laid-Open No. hei 06-131371 U.S. Ser. No. 08/131,943, which is configured to write music information in semiconductor memory used as a recording medium to enable reproduction of the music information any time. However, its interior circuit arrangement comprises an exclusive IC and an exclusive architecture, and as in the case of MDs, cannot realize wider use and extensive use of the device."}
{"text": "The present invention relates to a focus-detecting method and device. More particularly, the present invention relates to a focus-detecting method and device used in a microscope having reflected illumination.\nA conventional focus-detecting device used in a microscope with reflected illumination normally employs a detecting method in which a slit pattern is projected onto an object and a reflected image therefrom is detected by a photoelectric detecting device. Such a conventional arrangement is shown, for example, in U.S. Pat. No. 3,721,827, which is described hereinafter referring to FIGS. 9 and 10.\nFIG. 9 shows an example with a focal objective lens. Illumination light from a light source 1 passes through a condenser 2 and illuminates a slit 3. One-half of the beam emitted from the slit 3 is blocked by an illumination-side shielding plate 4 which blocks the beam from 0 degrees to 180 degrees around the optical axis. The other half of the light is reflected by a half mirror 6 and passes through an objective lens 7 to reach an object 8. The light reflected from the object 8 passes back through the objective lens 7 and is transmitted through the half mirror 6 to reach a photodetecting device 14.\nFIG. 10 shows a prior art example with an afocal objective lens. Illumination light from the light source 1 illuminates the slit 3 through the condenser 2. One-half of the beam emitted from the slit 3 is blocked by the illumination-side shielding plate 4. The other half of the light beam passes through a collimating lens 5 to form a parallel beam, is reflected by the half mirror 6, and passes through the objective lens 7' to reach the object 8. The light reflected from the object 8 passes back through the objective lens 7', is transmitted through the half mirror 6, and passes through a second objective lens 9 to reach the photodetecting device 14.\nBoth conventional examples shown in FIGS. 9 and 10 detect a focusing state by sensing the deviation of the center of gravity of the slit image formed on the photodetecting device 14. Not illustrated in FIGS. 9 and 10 are the eyepiece arrangement for viewing the object and the servomotor arrangement for automatically moving the object into the focal plane of the objective lens as a result of the image formed on the photodetecting device. These are conventional as shown in the aforementioned patent and are not illustrated in order to avoid unnecessarily complicating the drawings.\nHowever, the above-mentioned conventional technology has the following problem. When deviations in light intensity and phase distribution of an object 8 are large, as seen in a wafer pattern, both diffracted light and scattered light from the object 8 reach the photodetecting device 14. This results in poor focal accuracy. Generally focal accuracy is excellent when the light intensity and phase distribution of an object 8 are small as seen in an object having a mirror surface. The light shown by the dotted lines in FIGS. 9 and 10 is diffracted light and scattered light.\nTo solve this problem, JP Kokai H6-3578 discloses a focus-detecting device in which a section for forming parallel rays is provided after forming an intermediate image between the objective lens and the photodetecting device. A detection-side shielding plate is formed asymmetrically around the optical axis and placed therein.\nAccording to the technology disclosed in JP Kokai H6-3578, the diffracted light and scattered light which hinder focus detection are blocked. However, certain positioning of the detection-side shielding plate may block light which is necessary for focus detection. This makes it impossible to properly detect focusing."}
{"text": "1. Field of the Invention\nThis invention relates to methods and apparatus for eliminating clock jitter in Delta-Sigma, Analog-to-Digital Converters, and more particularly methods and apparatus for eliminating clock jitter in Continuous-Time Delta-Sigma Analog-to-Digital Converters.\n2. Description of Related Art\nRecently there has been a rapid rise in the use of Delta-Sigma analog-to-digital converters. They are capable of high-resolution data conversion without need for tight component matching or trimming of components. Because Delta-Sigma A/D converters are insensitive to a wide range of component non-idealities present in situations where they are implemented, high resolution sampling with relaxed component specifications is obtained by employing a single-bit digital-to-analog (DAC) inside a noise-shaping loop. The noise-shaping loop operates at a sampling rate that is typically much higher than the Nyquist rate of the data that is sampled. This causes sampling and component noise to be pushed to bands that are outside the band of interest. A paper by C. Candy and G. C. Temes, entitled \"Oversampling methods for A/D and DIA conversion,\" in Oversampling Delta-Sigma Data Converters, and published by IEEE Press, 1991, on pages 1-25 describes such converters and is incorporated by reference for its teachings on such. Due to the oversampling performed by Delta-Sigma converters, they typically require a simple anti-aliasing filter at their input.\nMany Delta-Sigma converters available and in commercial use today are Switched-Capacitor (\"SC\") implementations. Delta-Sigma converters are commonly employed as modulators. In the simplest case, a modulator may comprise a single Delta-Sigma converter. SC converter designs are used because the design and simulation of SC filters is well understood. Consequently, SC filters may be easily integrated into a design. However, because the capacitors employed in SC Delta-Sigma converter loops sample continuous valued voltage levels SC implementations require highly linear floating capacitors. In addition, both leads of these SC capacitors are generally switched during operation. These characteristics of SC designs complicate the manufacture of SC converters because the need for high linearity floating capacitors usually necessitates a double-poly manufacturing process.\nIn order to overcome these shortcomings, a new methodology called Switch-Current (\"SI\") has emerged as an alternative to SC designs for implementing sampled analog converter designs. SI designs are described in the following references: a paper by T. S. Fiez and D. S. Allstot, entitled \"CMOS Switched-Current Ladder Filters,\" published in IEEE J. of Solid-State Circuits, Vol. 25, No.6. December 1990, on pages 1360-1366; a paper by J. B. Hughes, N. C. Bird, and I. C. Macbeth, entitled \"Switched Currents: A New Technique for Analog Sampled-Data Signal Processing,\" published in Proceedings of the IEEE International Symposium on Circuits and Systems, May 1989, on pages 1584-1587; and a paper by J. B. Hughes and K. W. Moulding, entitled \"S-I: A Two-Step approach to Switched-Currents,\" published in the Proceedings of IEEE International Symposium on Circuits and Systems, 1993 on pages 1235-1238. These papers are hereby incorporated by reference for their teachings on SI designs. SI implementations advantageously can be at reduced implemented cost using standard CMOS processes. However, SI implementations have not yet achieved the same performance levels that have been achieved by the SC designs. Both SC and SI circuit implementations are examples of sampled-data systems. In a sampled-data system, a number of analog switches toggle at the rising edge of a clock signal. The sampled-data circuit is then given sufficient time to settle sampled voltages according to the circuit's switched topology. As a consequence, all voltages and currents are static in-between clock pulses in a sampled-data system.\nAn alternative to sampled-data systems is \"Continuous-Time\" (\"CT\") systems. In CT systems, sampling is performed at the converter output. As a consequence, the circuit nodes in the system operate in a CT mode. The implementation of Delta-Sigma modulators using the CT design approach advantageously performs sampling operations inside the modulator loop. Accordingly, sampling errors and out-of-band signals that alias into the passband are suppressed by the high gain of the loop in the passband. This advantageous effect is described in a paper by J. Schreier and B. Zhang, entitled \"Delta-Sigma Modulators Employing Continuous-Time Circuitry,\" published in EEE Transactions on Circuits and Systems-I: Fundamental Theory and Applications,\" Vol. 43, No.4, April 1996, on pages 324-332, which is incorporated by reference for its teachings on CT Delta-Sigma Modulators.\nThe above described effect is particularly important for passband or bandpass Delta-Sigma modulators where the input signal is varying at a rate close to or higher than the sampling rate. Bandpass modulators have become more frequently employed in IF sampling receivers. Further, because of their continuous nature of operation, CT circuits advantageously use relatively low power bandwidth amplifiers to achieve the same performance as SC and SI circuits. As a consequence, CT Delta-Sigma modulator circuits may have lower power consumption due to relatively low amplifier bias currents.\nDespite these advantages, two obstacles have impeded the use of CT circuits in Delta-Sigma modulator implementations. The first obstacle is that CT circuits are much more difficult to design and simulate than are SC and SI circuits. The second obstacle is that CT circuits are orders of magnitude more susceptible to clock jitter than are SC and SI circuits. With regard to the first obstacle, recently there has been significant research into the development of tools that map simulation results for sampled-data systems (such as SC & SI circuits) to their CT counterparts. In addition, CT simulations, although tedious and time consuming, are routinely performed by designers of ICs that implement CT RF, or base band filters. Therefore, the first obstacle is not an absolute deterrent against the use of CT Delta-Sigma designs.\nThe second obstacle, the susceptibility to clock jitter, is much more critical and has not yet been traversed by previous designs. The CT circuit's sensitivity to clock jitter has been a serious barrier against CT implementation of Delta-Sigma converters employed in modulators as noted in the following references: a paper by O. Shoaei and M. Snelgrove, entitled \"Optimal (Bandpass) Continuous-Time Sigma-Delta Modulator.\" published in the Proceedings of the IEEE international Symposium on Circuits and Systems, Vol. 5, May 30-Jun. 2, 1994, on pages 489-492; the Ph.d Thesis of L. Williams for the University of California, Berkeley, 1993 on pages 71-74; and in the Masters Thesis of N. Wongkomet, entitled \"A Comparison of Continuous-Time and Discrete-Time Sigma-Delta Modulators,\" for the University of California on pages 13-17. These references are incorporated by reference for their teachings on CT circuits and implementations.\nIn contrast to CT implementations, in sampled-data implementations clock jitter affects only the sampling capacitors at the front-end of the modulator. Once signals are sampled by these capacitors, clock jitter has no further effect in the sampled-data systems. This is due in part because the clock period in sampled-data systems is much longer than the required settling time. Hence, clock jitter effects for sampled-data systems may be modeled as modulation of input signals due to nonuniform sampling. Assuming the effects of clock jitter to be uncorrelated with the input signal and that the effects have a white noise spectrum, the clock jitter in a sampled-data system generates a white noise component that is added to the sampled signal. Due to an oversampling ratio (\"OSR\"), however, the noise caused by the clock jitter in the band of interest is attenuated.\nIn contrast, in CT implementations, clock jitter substantially adversely affects the operation of CT circuits. In contrast to the sampled-data designs, in CT circuits, the feedback DAC is implemented as a current source that is gated by the clock. In most Delta-Sigma modulator topologies, the signal from the DAC is subtracted from the input signal and thereafter a first integrator in the circuit integrates the result. As a consequence, in CT implementations, clock jitter generates a non-uniform ON-Time in the current source. The current source is directly summed with the input signal and is therefore, indistinguishable from the input signal. Assuming a spectrally \"white\" clock jitter, the noise added by the current source is proportional to the OSR because during every clock cycle an uncorrelated amount is added to the input signal. Therefore, the clock jitter tolerances for CT systems must be better than that of sampled-data systems by approximately the square of the OSR (e.g. -1000 for a typical OSR of 32). This results in clock jitter tolerances and specifications that are not easily achievable in most systems. Thus, a need exists for an apparatus and method that reduces or eliminates the deleterious effects of clock jitter on CT Delta-Sigma converters and the modulators employing the converters. In particular, a need exists for an apparatus and method that reduces the effects of clock jitter in CT implementations to levels comparable to those in sampled-data implementations. The present invention provides such an apparatus and methods."}
{"text": "This invention relates to a method for the removal of heavy metals from soil, mud, sludge and the like by the electro-kinetic process, whereby electric current is used to drive the heavy metals towards electrodes where they can be concentrated in an electrolyte solution from which they can be removed, and an apparatus for the execution of said method and the use hereof.\nThe need for cleaning soil and areas which are polluted with heavy metals is great and increasing in step with the recognition of the risks which these contaminated areas constitute in both the short and the long view for the environment.\nCleaning, however, is a very complicated and costly process. This means that in the most threatened areas the contaminated earth or sludge must be removed and deposited under controlled conditions for the safety of the surrounding environment.\nIn other words, the polluted material is isolated and thus the problem of cleaning is postponed until later.\nConsequently, there is a great need for a method of cleaning which is suitable for the removal of heavy metals.\nEfforts have been made with a so-called electro-kinetic method, whereby electric current is used to drive the heavy metals which exist in the earth in the form of ions towards electrodes where the metal can be concentrated and thereafter removed.\nIn practice, two electrodes are placed in separate chambers at a suitable distance in the earth. The positive heavy metal ions hereby move towards the negative electrode and the negative heavy metal ions towards the positive electrode.\nHowever, there is a difference in the mobility or the movement of the ions which stem from the electrode process and the heavy metal ions which are bound in the earth. The ions from the electrodes process are the least bound and are herewith of the easiest movability.\nThere thus arises the problem that a considerable part of the process current will be used to convey ions from the one electrode chamber to the other, so that only a smaller part of the power is utilized in the removal of the heavy metals in the earth.\nThis means that this known method is not particularly effective and, moreover, uneconomic in relation to the costs involved in the removal and the depositing of earth contaminated with heavy metals, which in turn means that the method has not been taken into practical use.\nAccording to the invention, by using a method whereby the earth is separated from electrodes in chambers, and that there is placed a cation-exchange membrane to prevent anions from the negative electrode chamber from seeping out into the earth and at the same time permit the passage of positive heavy metal ions from the earth and out in the negative electrode chamber, and also placing a second ion exchange membrane to prevent the passage of cations from the positive electrode chamber into the earth and at the same time permit the passage of heavy metal ions from the earth to seep out in the positive electrode chamber, there is achieved the advantage that the ions from the electrode chamber do not pass out into the earth, but remain in the electrolyte solution. The energy is hereby utilized in a far more efficient manner, i.e., solely from the transport of the earth-bound ions to the electrodes, and not for the transport of ions between the electrodes.\nAccording to the invention, it hereby becomes possible to clean products contaminated with heavy metals with a consumption of energy which is particularly competitive with other forms of cleaning, and at the same time to ensure a considerable degree of cleaning effect and herewith the possibility of achieving a degree of cleaning which is surprisingly high.\nBy using the inventive apparatus comprising an electrode in an electrolyte solution in a chamber separated from a second electrolyte solution in a second chamber by means of ion-exchange membranes, and separating this second electrode chamber from the earth by means of other ion-exchange membranes, an electrode unit which functions satisfactorily is achieved in a surprisingly simple manner.\nAn optimum effect of the electrode is ensured by isolating the electrolyte solution in the electrode chamber from the remaining electrolyte solution.\nBy using a circulating electrolyte solution, this can be pumped to a separate electrode chamber where the heavy metals can be deposited and removed from the electrolyte solution.\nFinally, it is expedient to use the method and the apparatus for the cleaning of earth, sludge and the like, either alone and/or in connection with a biological cleaning of the earth."}
{"text": "Cotton harvesters having on-board module forming structure such as described in commonly assigned U.S. Pat. No. 6,421,996 provide a compact bale or module directly on the harvester to reduce the amount of support equipment needed in the field and minimize harvester idle time during offloading. The on-board processing structure includes a bale handling system for moving a large formed bale rearwardly to prepare the bale chamber for a second bale with little or no harvester down time. The formed bale is relatively large, and designing the bale-forming chamber to maintain a satisfactory harvester shipping and transport height has been a continuing source of difficulty. The system shown in the aforementioned U.S. Pat. No. 6,421,996 patent provides reduced height by pivoting one portion of the bale chamber relative to the other portion so that the uppermost extremity of the chamber is lowered. The aft chamber can be supported on a rear bale handler and lowered for increased height reduction. As the capacity of a harvester increases, the problems with maintaining shipping and road transport heights within acceptable ranges become more acute. In addition, servicing the bale chamber area or replenishing the wrap mechanism of the baler, particularly in a chassis-mounted baler, becomes more challenging because of the increased height of the unit above the ground.\nWith previously available baling systems having pivoting gate structure to release a formed bale, separate latching mechanisms are required to lock the gate in place. Such mechanisms, which may include cylinder operated locks or complicated linkages which have to be adjusted and which add to the cost and complexity of the baling system."}
{"text": "A thiazole compound, typically, 2-chloro-5-(aminomethyl)thiazole, of formula (3):\nwherein X1 represents a hydrogen atom or a halogen atom, is a useful compound as an intermediate for pharmaceuticals and agrochemicals (see, e.g., JP 7-14916 B). For preparing the thiazole compound, several processes have been known. For example, (a) a compound of formula (1):\nwherein X1 is as defined above, and X2 represents a halogen atom, is reacted with hexamethylenetetramine, followed by hydrolysis (see, e.g., JP 4-234864 A and JP 4-21674 A); (b) the compound of formula (1) is reacted with potassium phthalimide, followed by hydrazinolysis (see, e.g., JP 4-234864 A); (c) the compound of formula (1) is reacted with formamide, followed by hydrolysis (see, e.g., JP 5-286936 A); and (d) the compound of formula (1) is reacted with ammonia (see, e.g., JP 4-234864 A and JP 2000-143648 A).\nHowever, the above processes (a) to (c) are not necessarily satisfactory from an industrial point of view because yields of the target thiazole compound of formula (3) are low in these processes. Although the above process (d) has an advantage over the processes (a) to (c) because of the use of low-priced ammonia, a considerable amount of a by-product compound of formula (4):\nwherein X1 is as defined above, is formed even when using ammonia in 20 mol-fold amount or more relative to the compound of formula (1), which results in a low yield of the target thiazole compound of formula (3). Therefore, further improvement has been desired."}
{"text": "1. Field of the Invention\nThe present invention generally relates to floating point units and, more particularly, to a floating point unit having an alternate, or fast, mode of operation which reduces pipeline stalls by issuing instructions regardless of the state of the previously issued instruction.\n2. Description of Related Art\nPipelining is a commonly practiced technique by which processing speed is increased by starting the execution of a next instruction sequence before a previous instruction sequence is completed. As the instructions within a pipe are at various stages of execution, correcting an error, commonly referred to as an exception, which occurs during the execution of a pipelined instruction can prove complicated. Accordingly, devices which use piplining are typically equipped to resolve exceptions in a manner which will minimize disruption to instructions flowing through the pipe.\nAn instruction execution pipeline of a processor typically includes instruction fetch, instruction decode, address calculation, execution and writeback stages. While both integer and floating point instructions are initially handled by the instruction execution pipeline, after completing the address calculation stage, floating point instructions are issued to a floating point unit (FPU) for the execution and writeback stages.\nIf the instruction execution pipeline has a floating point instruction in the pipe, the instruction will not be issued to the FPU until the previously issued floating point instruction has been checked for error. By delaying issuance of the floating point instruction, in the event that an exception occurs during execution of the previous floating point instruction, the FPU will be able to precisely identify and resolve the exception.\nWhile such a technique greatly simplifies the task of repairing the instruction execution pipeline, it often causes stalls in the flow of instructions through the pipe. More specifically, the address calculation stage will not issue a floating point instruction to the FPU until it is determined that the prior floating point instruction will not generate an exception. Thus, the address calculation stage waits to issue a next floating point instruction until after execution of the prior floating point instruction is complete. However, execution of most floating point instructions require more than one clock cycle. As a result, therefore, the floating point instruction stalls in the address calculation stage while waiting for the prior floating point instruction to finish.\nThus, it can be readily seen from the foregoing that it would be desirable to provide a method of issuing floating point instructions in a manner which reduces stalling of the instruction execution pipeline and a processor uniquely configured to permit operation of the FPU in plural modes. It is, therefore, the object of this invention to provide such a method and processor."}
{"text": "In a portable telephone system, in order to constantly ensure an optimum communication environment for wireless communication terminal that continues communication while its position is moved, a process for switching a base station (hereinafter, referred to as a handoff) to be connected is performed depending on a radio wave condition with the wireless communication terminal.\nThe wireless communication terminal measures both strength of a signal transmitted from the base station under communication and strength of a signal transmitted from other peripheral base stations other than that base station to carry out a process for maintaining more suitable radio state.\nFor instance, in the case of a CDMA portable telephone of an IS95 system (CDMA2000 1x system), the portable telephone measures a signal strength from a base station with which the portable telephone currently communicates and a signal strength from peripheral base stations other than that base station.\nThen, the portable telephone informs the base station of the measured signal strength while the portable telephone is under communication, and entrusts the base station to determine the handoff.\nFurther, during the idle state of the portable telephone, the terminal itself determines to switch the base station based on the measured result (for instance, see Patent Document 1).    Patent Document 1: JP-A-2002-171555"}
{"text": "Some distributed computing environments, such as Parallel Sysplexes, today provide a non-volatile shared storage device called the coupling facility, that includes multiple storage structures of either the cache or list type. These structures provide unique functions for the operating system and middleware products employed for the efficient operation of a Parallel Sysplex. For example, the cache structures provide directory structures and cross-invalidation mechanisms to maintain buffer coherency for multisystem databases, as well as a fast write medium for database updates. These are used by, for instance, the data sharing versions of DB2 and IMS, offered by International Business Machines Corporation, Armonk, N.Y.\nThe list structures provide many diverse functions. One such list structure function is to provide for high-performance global locking, and this function is exploited by such products as the IMS Resource Lock Manager (IRLM) and the Global Resource Serialization (GRS) function in OS/390, offered by International Business Machines Corporation, Armonk, N.Y. Another list structure function is to provide a message passing mechanism with storage for maintaining multiple messages on a per system basis and a mechanism for notifying a system of the arrival of new messages. This function is exploited by the XCF component of OS/390, which in turn is exploited by numerous multisystem applications for providing a capability to pass messages between their various instances. A third list structure function is to provide for shared queue structures that can be ordered and accessed by LIFO/FIFO ordering, by key, or by name. Workload Manager (WLM), IMS Shared Message Queues and MQ Series, all offered by International Business Machines Corporation, Armonk, N.Y., are examples of exploiters of this feature. While these functions provide examples of the list structure uses, other uses exist.\nVarious components of a Parallel Sysplex have been documented in numerous applications/patents, which are listed above and hereby incorporated herein by reference in their entirety. The capabilities defined in some of those patents provide the basic system structure to create and manage cache and list structure instances. Additionally, various of the applications/patents listed above provide extensions to the base functions of the Parallel Sysplex.\nIn many situations, a failure of the coupling facility that contains various structures requires significant recovery actions to be taken by the owning applications. For example, for database caches and queues, this may require using backup log data sets and/or tapes. This is a time-consuming process that results in a loss of access to the application during the recovery operation. Other structures, such as lock tables, may require reconstruction of partial lock tables from in-storage copies, along with failures of in-flight transactions. Still other structures, such as message-passing structures, may lose all data and require re-entry from the application. So, there is a proliferation of diverse recovery schemes with different recovery times and impacts. Moreover, since the failure of a coupling facility results in all resident structures failing, the diverse recovery actions are occurring concurrently, which can cause serious disruptions in the Parallel Sysplex.\nThus, a need exists for a configuration of a Parallel Sysplex that provides less disruptions. In particular, a need exists for a high-availability coupling facility, which improves on the recovery times and impacts of existing recovery techniques, while also provides for a consistent recovery design across various structure types. As a particular example, a need exists for one or more capabilities that facilitate duplexing of structures in separate coupling facilities coupled to one another."}
{"text": "A substrate processing apparatus is used in the process of manufacturing a semiconductor, a liquid crystal panel, or the like, to treat a surface to be treated of a substrate such as a wafer or a liquid crystal substrate with a chemical solution, rinse the surface after the chemical treatment with a rinsing liquid, and then dry the substrate after the rinsing.\nIn the drying process, problems such as pattern collapse around, for example, the memory cell or the gate have occurred due to recent miniaturization along with higher integration and higher capacity of semiconductors. This is caused by the spacing and structure between the patterns, the surface tension of the rinsing liquid, and the like.\nTherefore, in order to suppress the pattern collapse, there has been proposed a substrate drying method using a volatile solvent (e.g., IPA: 2-propanol, isopropyl alcohol) with a surface tension smaller than that of the rinsing liquid (e.g., DIW: ultrapure water). In the substrate drying method, the substrate is dried by replacing the rinsing liquid on the treated surface of the substrate with a volatile solvent. At this time, to promote the drying, the substrate may be heated by a lamp located above the surface to be treated of the substrate to dry the substrate.\nIn the chemical treatment step, a high-temperature treatment of 150° C. or more may be performed using sulfuric acid or phosphoric acid as a chemical solution. In this case, the chemical solution is heated to 150° C. or higher to be used in the treatment. Therefore, the chemical solution on the surface to be treated of the substrate is heated to a high temperature of 150° C. or higher by the lamp located above the surface to be treated.\nHowever, in the drying process and the chemical treatment process, since the lamp is located above the surface to be treated of the substrate, a mist or particles of the treatment liquid adhere(s) to the surface of the lamp during the substrate processing. The adhering matter such as a mist or particles of the treatment liquid drops from the surface of the lamp and adheres to the surface to be treated of the substrate. This may result in the contamination of the substrate.\nMeanwhile, a common substrate processing apparatus is provided with an ultra-low penetration air (ULPA) filter or a high efficiency particulate air (HEPA) filter. The clean air that has passed through the filter flows in the substrate processing apparatus as downflow (vertical laminar flow), and the inside of the substrate processing apparatus is kept clean. As a result, particles are prevented from adhering to the inside of the substrate processing apparatus, in particular, the surface to be treated of the substrate.\nHowever, in the aforementioned substrate processing apparatus, since the lamp is located above the surface to be treated of the substrate, clean air directed from the filter toward the surface to be treated of the substrate is interrupted by the lamp and is not sufficiently supplied to the surface to be treated of the substrate. For this reason, it is difficult to remove particles from the surface to be treated of the substrate and its surroundings. As a result, the substrate is contaminated."}
{"text": "1. Field\nSubject matter disclosed herein relates to techniques for reporting on or tracking vehicles associated with a particular area. In some implementations, reporting and tracking techniques are provided for ground vehicles in an airport environment implementing Automatic Dependent Surveillance-Broadcast (ADS-B) technology.\n2. Information\nIt may be desirable to have current information regarding ground vehicles that are within and/or near a particular area. One type of area in which information may be desirable comprises an airport. For example, security concerns or airport operational issues may make it beneficial to know what type of vehicles may currently be in an airport, where vehicles may be located, how fast they may be moving, what direction they may be moving, who may own the vehicles, and/or other information. Techniques are needed for accurately tracking vehicles within a specific region."}
{"text": "The invention relates to a communication system and method. The system involves a novel transmitting system and method and a novel receiving system and method.\nData compression and transmission is highly desirable as saving bandwidth or time of transmission or both. Accordingly, efforts have been made to reduce the time of transmission and compress bandwidth by various means.\nU.S. Pat. No. 3,334,335, issued Aug. 1, 1967 to Brick et al. for \"Electronic Data Processing\" mentions a system in which a word is received in one language and produces an address as an output in response to the word input. The patent itself is concerned with a system in which a second or foreign language word is read out. It also describes systems in which a memory stores each word, an input word being compared with the stored word and a signal generated when a match is found. This match signal is subsequently used to determine the address of the correct stored word. The patent itself is concerned with the indication for match and mismatch of the words.\nIn another patent, U.S. Pat. No. 3,717,851, issued Feb. 20, 1973 to Cocke et al. for \"Processing of Compacted Data,\" compacted data in the form of variable length codes are employed. These have length-representing prefix portions which are variable-length encoded. The variable length prefixes are decoded by small, fast, search-only associative memories which provide a match-indicating signal as an address to a second memory having conventional storage elements. The output of the last main memory may contain a base address to still another memory and a indication of the bits remaining in the variable length code word.\nIn U.S. Pat. No. 3,400,380, dated Sept. 3, 1968 to Packard et al. for \"Digital Computer Having an Address Controller Operation,\" there is described a computer system in which digit, or four bit data representations, are transferred to character or eight bit format.\nOther examples of conversion systems are described in U.S. Pat. No. 3,229,047 of Jan. 11, 1966 to Simpson for \"Data Conversion Systems\"; U.S. Pat. No. 3,662,347 of May 9, 1972 to Fox for \"Signal Compression and Expansion System Using a Memory\"; U.S. Pat. No. 3,778,617 of Dec. 11, 1973 to Calhoun for \"Container Inspection Apparatus\"; and U.S. Pat. No. 4,032,913, issued June 28, 1977 to Okamura et al. for \"Coding Equipment Providing Compressed Code.\""}
{"text": "With impetus on going green, energy mix is changing rapidly. Today, the energy mix is getting a lot more volatile given the very nature of such energy sources like Wind and Solar as they are exposed to prevailing weather condition. Unlike traditional sources such as thermal generation systems which typically have spinning reserve; all such integrated and/or synchronized energy sources contribute to the overall grid inertia. With depleting grid inertia, any imbalance in power position either due to demand side or the supply side creates havoc to the grid frequency. It may be noted that the demand of energy has always been volatile as they are exposed to exogenous factors such as weather changes. With supply of energy also getting volatile, maintaining the grid health is becoming more and more challenging and costlier day by day.\nThe situation may become much more challenging as in today's world where consumers are themselves become energy generators, it becomes much more difficult to monitor and manage such a vast network having enormous number of generation units. While smart meters and/or smart grids are enabled to capture such data in real time and gives two-way control on the grid, there is a growing need for real time intelligence to be drawn from such data for more proactive and efficient decision making by the energy utilities across the value chain i.e. generation, distribution and transmission."}
{"text": "The present invention relates to a memory circuit, and more particularly to a dynamic memory circuit having a plurality of memory cell arrays.\nThe storage capacity of semiconductor memory devices has been becoming greater and greater, and efforts are being made to increase the integration density of integrated circuit chip.\nHowever, if such increase in the storage capacity is achieved by simply increasing the number of memory cells coupled to each word line, the storage capacitance of each word line would become large, resulting in low speed operation.\nFurthermore, in the case where the memories are of dynamic type, the written information disappears when a predetermined time has elapsed. Therefore, it is necessary to execute a rewriting cycle called \"refresh\" within the above time. The refresh is done on the respective word lines one by one and hence, the number of the word lines in the memory corresponds to the cycle times of refresh within the predetermined period of time.\nUnder the above circumstance, it has been practiced that the memory cells are divided into two or more memory cell arrays each having word lines to which reduced number of memory cells are coupled. In such construction, refresh operations are executed in the respective memory cell arrays in parallel, reducing the time required for refresh for the whole memory cells.\nFor example, memory cells of a 64 K-bit dynamic RAM are arrayed in two memory cell arrays each having 128 word lines and 256 bit lines, and the whole memory cells are refreshed through 128 refresh cycles. In this kind of memories, one word line in each memory cell array is driven and two word lines are simultaneously driven as a whole in the memory in response to the same driving pulse. Accordingly, the drive circuit for generating the driving pulse is required to have a large driving capacity in order to attain high speed operation. Moreover, since two word lines are driven at the same time to be charged to a power voltage, a large value of peak current necessarily flows to the two word lines through the drive circuit. This peak current causes noise and affects the memory operation, resulting in malfunction."}
{"text": "A digital television (DTV) is now presented to offer various services in addition to a television (TV)'s original function such as playing video and audio. For example, broadcasting information such as Electronic Program Guide (EPG) may be provided to a user, and also, broadcasting services from at least two channels may be simultaneously provided to a user. Especially, since a receiving system of the DTV includes a large capacity of a storage device, and is connected to a data communication channel and the internet (through which two-way communication is available), more services become accessible through broadcast signals. Additionally, since services offered through broadcast signals become more diversified, needs for utilizing the diversified services accurately is increased.\nRather, since an image display device in each home is connected to a broadcast receiving device such as a settop box, it may play an uncompressed AV content that the broadcast receiving device provides. A broadcast receiving device receives content from a server called a multichannel video programming distributor (MVPD), and after extracting content from a broadcast signal received from the broadcast receiving device, the MVPD converts the extracted content into a signal of a format fit for transmission and provides the converted signal to the broadcast receiving device. In such a way, since an image display device less likely receives a broadcast signal directly through airwaves, a broadcasting station transmitting a broadcast signal requires a method of providing self-led enhanced service to the image display device."}
{"text": "1. Field of the Invention\nThe present invention relates generally to resource sets. Still more particularly, the present invention changes the way resource sets are represented.\n2. Description of the Related Art\nA resource set (RSET) is a structure that identifies physical resources. An RSET parameter issued in resource set application protocol interfaces (APIs) or commands is to either get information from the system regarding resources or to pass information about requested resources to the system. Applications and job schedulers can attach an RSET to a process. Attaching an RSET to a process limits the process to use only the resources contained in the RSET. The central processing unit (CPU) and memory resources in a resource set are represented by bitmaps. The primary use of RSETs is to perform CPU topology and affinity operations. CPUs are identified in RSETs by logical CPU IDs. For example, a partition with two multi-chip modules (MCMs) of eight processors each can have their 16 logical CPU IDs assigned in any order. Using RSET APIs, applications cannot assume that logical CPU IDs 0-7 are in one MCM and IDs 8-15 are contained in the other MCM.\nThe set of logical CPU IDs available in a system might not be contiguous. There can be gaps in logical CPU ID numbers. This can occur when CPUs are dynamically reconfigured out of a partition. Some systems allocate logical CPU IDs for the online CPUs sequentially at boot time. The main system-defined resource sets are:                System RSET and sys/sys: This resource set contains the available (online) CPU and memory pool resources in the system or partition. On partitionable machines, this RSET contains only the resources that are in the operating system's partition.        Node RSETs, sys/node.mm.nnnnn, or sys/node.nnnnn: These RSETs contain resources that are present at various system detail levels (mm) and indexes (nnnnn) in the system. For example, if system detail level 04 represents the level in the system topology that corresponds to a p690 MCM, RSET sys/node.04.00000 contains the resources in an MCM 0. RSET sys/node.04.00001 contains the resources in another MCM 1, and so on.        Atomic RSETS, sys/cpu.nnnnn, or sys/mem.nnnnn: These RSETs contain a single resource, either a CPU or memory pool. There are atomic resource sets for every available (online) resource contained in the operating system's partition.        \nThe resource set configuration is displayed with the lsRSET command. There are two types of RSET, the partition RSET and the effective RSET:                The partition RSET can only be attached, modified, or detached by a root user. Some managers attach a partition RSET when a process is classified with a work class that contains an RSET. There is only one partition RSET per process and it is updated by replacement. For example, a process is started with a specific class that attaches a partition RSET that contains CPUs 0-3. Later a root user attaches an RSET that contains CPUs 2-7. The partition RSET attached by manager is replaced by the new RSET. The process now runs on CPUs 2-7.        The effective RSET, generally used by applications, can be attached by root users and non-root users with a CAP_NUMA_ATTACH. Effective RSET limits a process to run only on the resources (CPUs, memory) contained in the RSET. This means that the effective RSET of a process cannot contain more resources than the process' partition RSET. For example, a process may have a partition RSET established by the WLM that limits the process to running only on CPUs 0-3. A user can attach an effective RSET with CPUs 2-3, and the process is limited to running only on CPUs 2-3. An attempt by the user to attach an effective RSET with CPUs 2-7 would be rejected because the user attempted to use resources outside its partition RSET.        \nOnly recently has the partition RSET been used. In the future, job schedulers may also use partition RSET. The RSET commands provide an easy way for system administrators to use system RSETs. Commands are provided to make, display, and remove RSETs from the system registry. Other commands allow RSETs to be attached to running processes or a command attached to an RSET.\nThus, the concept of an RSET is used to represent a subset of the system's physical resources such as CPUs. Once an RSET is created, it can be associated with logical resources such as a virtual memory object or a process. The operating system than restricts those logical resources to the physical resources contained within the RSET to which they are associated. For example, a process can use an RSET to attach itself to the CPU described in the RSET. Each element in the SRET describes a CPU, and the process attaches to all the elements.\nCurrently, the set of resources described by the RSET remains unchanged, or static, throughout the life of the RSET. However, in an environment where the system configuration is dynamically changing; processors and memory can be added and removed from the system. An RSET once created does not change to reflect changes in the system configuration. As an example of this problem is where using a RSET, a bind is made between a process and a set of CPUs that have an affinity with one another. A processor re-configuration operation can subsequently add another process that has an affinity to the CPUs in the processes RSET. But the process will remain bound only to those processors that are contained in the RSET and will not be dispatched on the additional processors that were added to the node dynamically. It may have been an administrator's intention to run the process on all CPUs that have the same affinity, but since the view changed, the information was lost to the process.\nThus, it would be advantageous to have an improved method, apparatus, and computer usable code that allows an RSET to be used in a dynamic environment."}
{"text": "Medical devices may be used to treat a variety of medical conditions. Medical electrical stimulation devices, for example, may deliver electrical stimulation therapy to a patient via implanted electrodes. Electrical stimulation therapy may include stimulation of nerve, muscle, or brain tissue, or other tissue within a patient. An electrical stimulation device may be fully implanted within the patient. For example, an electrical stimulation device may include an implantable electrical stimulation generator and one or more implantable leads carrying electrodes. Alternatively, the electrical stimulation device may comprise a leadless stimulator. In some cases, implantable electrodes may be coupled to an external electrical stimulation generator via one or more percutaneous leads or fully implanted leads.\nMedical electrical stimulators may be used to deliver electrical stimulation therapy to patients to relieve a variety of symptoms or conditions such as chronic pain, tremor, Parkinson's disease, depression, epilepsy, urinary or fecal incontinence, pelvic pain, sexual dysfunction, obesity, or gastroparesis. An electrical stimulator may be configured to deliver electrical stimulation therapy via leads that include electrodes implantable proximate to the spinal cord, pelvic nerves, gastrointestinal organs, peripheral nerves, or within the brain of a patient. Stimulation proximate the spinal cord and within the brain are often referred to as spinal cord stimulation (SCS) and deep brain stimulation (DBS), respectively.\nA clinician selects values for a number of programmable stimulation parameters in order to define the electrical stimulation therapy to be delivered to a patient. For example, the clinician may select a current or voltage amplitude of the stimulation, and various characteristics of the stimulation waveform. In addition, the clinician may specify an electrode configuration used to deliver stimulation, including selected electrode combinations and electrode polarities. If the stimulation is delivered in the form of pulses, for example, the clinician may specify a current or voltage pulse amplitude, pulse width and pulse rate. A set of parameter values may be referred to as a stimulation program. A program group may include multiple programs. Multiple programs in a program group may be delivered on a simultaneous, time-interleaved, or overlapping basis."}
{"text": "1. Field of the Invention\nThe present invention relates to data acquisition systems, and more particularly to a digitally configurable multiplexer/de-multiplexer for routing a number of analog signals to or from selected analog front ends of a digitization block of a data acquisition system.\n2. Description of the Related Art\nAnalog multiplexers/de-multiplexers (MUX/DEMUX) are typically configured as D:1 devices (e.g., 2:1, 4:1, 8:1 or 16:1) with “N” address inputs for selecting from among “D” analog data signals in which D=2N. An analog MUX/DEMUX is typically bidirectional and may be operated as a MUX to select from among multiple inputs to an output or as a DEMUX to route an analog signal to redirect to a selected one of multiple outputs as known to those skilled in the art. The term “multiplexer” or “MUX” as used herein is intended to connote bidirectional operation in which the MUX may be used for MUX or DEMUX operations. Each configuration has traditionally been implemented on a separate integrated circuit (IC) or chip. If an analog signal needed to be dynamically re-routed to another path, additional multiplexers or switches were required and were placed in series with the signal path. Such additional devices added cost and complexity to the overall data acquisition system. In many cases, different configurations are required in the same system, which forced the designer to implement and purchase different ICs that suited each of the configurations. For example, some systems required the same signal to be routed to any number of different analog front ends that performed various functions (e.g., filtering, amplifying, etc.) depending upon the characteristics of the signal.\nConventional analog multiplexers are relatively inflexible and consume a significant amount of printed circuit board (PCB) area of the data acquisition system."}
{"text": "The access point (AP) should be able to support a large number of associated stations (STAs) (e.g. over 6000) which may operate on a very strict energy budget. Such devices could be battery powered sensors which transmit and receive data very rarely and stay in the low power operation mode for relatively long periods of time. The basic operation mode is the distributed coordination function (DCF) mode. In order to support a large number of STAs operating in random access mode the access point may utilize specific techniques to restrict contention to the channel to avoid of collisions of simultaneous transmissions. One such technique to reduce contention and collisions of transmission is a method of grouping STAs to multiple groups and assigning certain parameters for each group to indicate the specific group that can access content for the channel. Such grouping information and parameters for the operation can be delivered to the STAs in an association phase or in broadcast information such as beacons.\nThe STAs may operate in low power mode for prolonged periods of time and therefore the grouping related parameters may be not be valid anymore when a STA wakes up and resumes the channel access operation. Broadcast messages such as beacons are transmitted only occasionally and if a STA stays in low power mode for extended periods of time it may lose synchronization (due to the internal clock drift in the AP and in the STA) and is not able to estimate when the next beacon would be transmitted. This may cause the STA to stay awake for relatively long period trying to receive a beacon transmission.\nIn 802.11 protocol the AP buffers data frames if the STA is in the low power state. The AP informs the STAs about the buffered frames by indicating buffered data in the traffic indication map (TIM) which is transmitted in the beacon message. Once the STA awakes it can retrieve the buffered frames by transmitting some buffered uplink (UL) data which implicitly notifies the AP that the STA is awake, or the STA can transmit a power save poll (PS-poll) to indicate to the AP that it is awake and ready to receive data. Alternatively any frame that is classified to be a so called ‘trigger frame’ indicates that STA is awake.\nTo improve data exchange between AP and STAs, speed frame exchange can be initiated between AP and STA (or between STA and another STA). AP may have data in the DL buffer but is not able to transmit all the data in the same session since it may have some pending management frames in the buffer (such as a beacon transmission) or other control/data frames it may need to deliver."}
{"text": "(1) Field of the Invention\nThe present invention relates to a skew compensation method and a skew compensation apparatus for automatically compensating for skew occurring among a plurality of serial transmission paths used for high-speed signal transmission with parallel-to-serial conversion.\n(2) Description of the Related Art\nParallel-to-serial transmission is a technology which serializes a parallel signal by time division multiplexing, transmits the resultant serial signal, and then parallelizes the serial signal again at a receiver. The technology allows the elimination of the problem of bit-to-bit skew which is conspicuous when the rate of parallel transmission is increased. To implement parallel-to-serial transmission in a broader band, it is necessary to increase the rate of serial transmission, but the acceleration of serial transmission is limited. For broader-band parallel-to-serial transmission, therefore, an approach of using a plurality of serial signals in parallel has been proposed. However, the approach entails the occurrence of skew between the serial signals, similarly to the parallel transmission.\nIt is inherently possible to suppress the skew by installing each of the plurality of transmission paths over an equal distance. Due to reduced margin resulting from higher-speed signal transmission, however, it is difficult to install all the lines for serial signals such that they have equal lengths. In particular, tremendous skew occurs in proportion to a transmission distance in the case of WDM (Wavelength Division Multiplexing) used for optical communication. In the WDM, a plurality of, serial signals are transmitted in parallel by utilizing the property of a single optical fiber which carries a plurality of optical signals with different wavelengths. In this case, however, the different wavelengths cause tremendous skew irrespective of equal transmission distances.\nJapanese Unexamined Patent Publication No. 2000-341135 discloses a technology which transmits a test pattern in equiphase from each of the transmitters for serial signals to eliminate skew occurring among a plurality of serial signal transmission paths. In accordance with the invention disclosed in the precedent document, each of receivers sends the timing of a synchronization signal to a controller. The controller sets the delay line of the transmitter with the most delayed one of the received timings and thereby matches the reception timings at the receivers.\nHowever, since the foregoing method sends a test pattern by switching normal data to the test pattern, the skew cannot be monitored during the transmission of the normal data so that a sudden skew shift is beyond handling."}
{"text": "The present invention relates to a positive displacement machine.\nPositive displacement machines have been widely used in practice. Such positive displacement machines as piston pumps or rotary pumps are composed of a great number of separate parts, and for this reason they are very expensive and bulky. Relatively simple gear pumps which are used when it is necessary to provide a high pressure, are even more expensive inasmuch as voluminous sealing elements are provided on the side surfaces of the gears. They are not always satisfactory in operation, and the construction costs of high pressure gear pumps are great."}
{"text": "A roller assembly, roller unit, and conveyor apparatus which convey a convey target must intermittently convey a convey target. For this purpose, there is provided a technique for adjusting an idle load using a frictional force set such that roller rotation is stopped when the load of the convey target is larger than the convey force of the roller and the roller is rotated to convey the convey target when the load of the convey target is smaller than the convey force of the roller.\nFor example, in PTL1, a sliding bearing is disposed between a driving roller fitted on a shaft portion and a convey roller. When a driving force is transmitted from a driving source to the driving roller, the sliding bearing is driven to rotate upon rotation of the driving roller by a frictional force between the sliding bearing and the driving roller. Upon rotation of the sliding bearing, the driving force is transmitted to the convey roller by the frictional force between the sliding bearing and the convey roller."}
{"text": "1. Field of the Invention\nThe present invention relates generally to fire fighting tools and, more particularly, to a fireman's axe. Specifically, the present invention relates to an improved combination fire tool performing the functions of a plurality of fire fighting tools previously required heretofore.\n2. Description of the Prior Art\nFire fighters are frequently confronted with unpredictable situations wherein a number of different types of tools might be required. Fire fighters commonly faced the problem of making quick entries into buildings or quick exits through walls or doors in emergency situations. Gaining entry into buildings often requires a prying action to dislodge members from doorways and windows. A hammering action is also required for breaking out windows and doors to gain entry. It is also frequently required to pull down ceilings of drywall or lathe and plaster in order to investigate or gain access to fire location in a building.\nAs a result of the above, most fire fighters find they may need considerably more than just one particular tool but must make a choice of carrying only one or two. Heretofore, several fire fighters needed to coordinate with each other to bring into the fire area an axe, a sledgehammer, a pry bar and/or a pike pole in order to accomplish any of the aforementioned tasks. Consequently, several fire fighters have needed to carry in a plurality of different tools in order to have the proper number and types of tools available. Unfortunately, tools that are not being used at a particular time must be set down and are frequently lost due to their being forgotten or covered by falling debris. As a result, numerous tools needed on site are lost during the course of fighting a fire.\nCombination firefighter tools have evolved-over the years in order to attempt to alleviate some of the duplication of separate tools discussed above. U.S. Pat. Nos. 3,219,316, 3,599,255, 3,604,028, 4,287,623, 4,597,123, 5,044,033, and Des. No. D120,609 and No. D233,405 all disclose various fireman combination tools. Unfortunately, such combination tools have generally been limited in their multiple purposes and have generally been in the size and form of hand axes. While the device disclosed in U.S. Pat. No. 4,287,623 has a plurality of different purposes, its complexity of parts and operation defeats the purpose of having one sturdy tool at a fire site.\nOther utility patents disclosing combination tools, generally in the form of hand axes, includes U.S. Pat. Nos. 89,013, 292,168, 637,253, 790,973, 1,596,602 and 4,030,150 as well as Des. Nos. D35,154, D45,761, 48,231, D67,749, D163,911 D299,414. These remaining patents, both design and utility, illustrate a variety of combination tools which also attempt to serve a multiplicity of purposes. However, none of the aforementioned references disclose a combination fire tool which serves multiple purposes so as to avoid the requirement of separate axes, sledgehammers and pike poles as well as other tools. Therefore, there remains a need for a combination fire tool which is sturdy and simple in design and function."}
{"text": "Improvements in computer processing power and broadband technology have led to the development of interactive systems for navigating imagery, such as geographic imagery. Some interactive map navigation systems provide a user interface with navigation controls for navigating cities, neighborhoods, and other terrain in two or three dimensions. The navigation controls enable users to tilt, pan, rotate, zoom, and activate terrain and buildings for different perspectives at an area of interest. Exemplary map imaging systems include the Google Earth and Google Maps systems by Google Inc.\nThe imagery used by these interactive systems can be derived by rending geometric objects and texture objects to generate a two-dimensional or three-dimensional graphical representation of an area of interest. The geometric objects can define three-dimensional surfaces of objects and terrain depicted in the imagery. The texture objects can be mapped to the three-dimensional surfaces to add detail, surface texture, color and other features to the depicted objects and terrain.\nInteractive systems for displaying imagery, such as geographic imagery, often render geometric and texture objects with low level of detail (“LOD”) (e.g. low resolution) at camera views far from an area of interest to improve performance. As the user zooms in closer to the area of interest, geometric and texture objects with higher LOD (e.g. higher resolution) are rendered to depict objects in greater detail. Switching between two levels of detail can cause a distracting pop if the levels of detail are sufficiently different and/or the transition happens too close to the camera view. The pop can be a visual anomaly which disrupts the user's sense of immersion in the imagery. Moreover, if the geometric data and the texture data are fetched from a server over a network, the latency between the request for higher resolution objects and the receipt of the data by the client can be long. In this case, a user can zoom in past the natural resolution of a rendered scene, resulting in a blurry image. When the higher resolution imagery finally arrives from the server and is rendered, the higher resolution imagery will “pop” into view because its resolution is markedly different from the imagery in the previously displayed blurry image."}
{"text": "There has long been motivation for alternative energy sources to reduce the dependency on non-renewable energy sources such as oil and coal. There has long been interest in harnessing wind to generate electricity particularly in areas with a high or relatively continuous stream of wind energy.\nThere have been challenges and previously proposed solutions for harnessing wind energy through use of electric wind turbines. These large windmill-type structures suffer from many disadvantages including, the ability to adjust to variable wind conditions in either or both of direction or velocity. Additional disadvantages of prior wind turbines include relatively low efficiency rates due to standard designs that are employed in environments not suitable for the particular wind turbine design or initial set-up. Due to the size and mass of large wind turbines, premature structural fatigue and failure of components occur. Prior wind turbines also suffer from inefficient and complex designs.\nThus it is desirable to improve on prior wind turbine designs to increase the turbine's ability to harness the wind energy, improve the efficiency in converting the wind energy to electricity and to simplify and improve their design, manufacturing and operation."}
{"text": "Authentication for computer devices can occur in many different ways. Still today, authentication using a username and password is most common. While this type of authentication still works, it is often complicated and cumbersome to enter a username and password which is secure enough to be difficult to break.\nUS-2010/218249 presents a system and/or a method that facilitates authentication of a user in a surface computing environment. A device or authentication object can be carried by a user and employed to retain authentication information. An authentication component can obtain the authentication information from the device and analyze the information to verify an identity of the user. A touch input component can ascertain if a touch input is authentication by associating touch input with the user. In addition, authentication information can be employed to establish a secure communications channel for transfer of user data. However, access control to computers always benefit from improved usability with the same or higher security.\nEP-2 395 446 presents a method for pairing a first device (such as a handset) with a second device (such as a personal computer), said first and second devices being linked to a remote server, said first device comprising an accelerometer and said second device being linked to a third device. The method comprises the following steps: the first device and the second device capture a common movement respectively to the first and to the third devices; the first device and the second device send data resulting from the common movement to the remote server so as to create a unique identifier of the movement; the remote server compares the identifier of the movement of the first device and the identifier of the movement of the second device, and if the identifiers match, the remote server pairs operations performed from the first and the second devices. However, this method requires the user to hold the first device while performing a movement using e.g. a mouse or a touchpad which can be correlated. Such a correlation is cumbersome and resource demanding."}
{"text": "1. Field of the Invention\nThe present invention relates to magnetic storage devices generally, and in particular to a crosstie random access memory element and circuitry, for reading and/or writing at a application Ser. No. 386,388 filed Jun. 8, 1982 with respect to which the present application is a continuation. This application is also related to U.S. Pat. No. 4,962,477 to Schwee, which in turn refers to other related copending applications, including Ser. No. 386,389 of Schwee et al filed Jun. 8, 1982 (abandoned in favor of its continuing application, Ser. No. 07/615,728, filed Nov. 19, 1990) naming Paul E. Hunter as coinventor who is not a coinventor of subject matter covered by the present application.\n2. Description of the Prior Art\nElectronic computers and other data processing devices, to a great extent, are limited by the speed, capacity, and reliability of their memory systems. Memory systems currently in use include tiny ring shaped ferrite cores strung on a mesh of fine wires (core memories) and semiconductor memories comprising transistor circuits laid down on tiny chips of silicon. At the present time the most common nonvolatile random access memory (RAM) in use is core memory. However, core memories having large storage capacity become bulky, heavy and extremely expensive compared to memories which can be fabricated using photolithographic techniques such as the aforementioned semiconductor memories. Unfortunately, for the most part, dense and light weight semiconductor memories are volatile. The exception is MNOS semiconductor memories which take a relative long time to write and read and have limited retentivity Consequently, there is a need in the prior art to configure by use of photolithographic techniques a nonvolatile random access memory in association with read/write circuitry, which is fast, but light in weight and low in cost.\nIn the prior art, as an alternative to the foregoing memories data bits are stored in the form of magnetic bubbles which move in thin films of magnetic material These magnetic bubbles are actually cylindrical magnetic domains whose polarization is opposite to that of the thin magnetic film in which they are embedded The magnetic bubbles are stable over a considerable range of conditions and can be moved from point to point at high velocity. Magnetic bubble memories are substantially cheaper than the core memories aforementioned, and much faster than magnetic disk memory systems which are widely used for high capacity storage. Magnetic bubble memories are analogous to magnetic disk memories in that both systems store data as states on, or in, a thin magnetic film. In a disk memory, the film is moved mechanically at high speeds while in a magnetic bubble memory, the magnetic bubbles move at high speeds throughout the thin magnetic film. Since magnetic bubble memories are nonvolatile, logical operations can be performed without reading or writing the stored data out/in again so that magnetic bubble memories do not have moving parts and therefore may operate reliably for long periods of time. However, a drawback of magnetic bubble memory devices is that single crystals with minimal defects are required for fabrication. Moreover, it is difficult to produce large single crystals for high storage devices. Even the attempted use of amorphous materials does not eliminate the problem of high eddy currents which reduce the speed of the magnetic bubble domains. Even so, magnetic bubble memories have many advantages. Nevertheless, there is a need in the prior art to configure a nonvolatile random access memory using thin magnetic film technology and yet provide rapid access time while utilizing only polycrystalline material.\nAs further background, the present invention is an out growth of the crosstie memory technology disclosed primarily in U.S. Pat. Nos.: 3,846,770 to Schwee et al., entitled \"Serial Access Memory Using Magnetic Domains in Thin Film Strips,\" filed Jul. 11, 1973, patented Nov. 5, 1974; 3,868,659 to Schwee, entitled \"Serial Access Memory Using Thin Magnetic Films,\" filed Apr. 10, 1973, patented Feb. 25, 1975; 3,868,660 to Schwee, entitled \"Detector for Cross-tie Memory,\" filed Apr. 10, 1973, patented Feb. 25, 1975; 4,100,609 to Schwee, et al., entitled \"Magnetoresistance Detector for Crosstie Memories,\" filed Sep. 3, 1976, patented Jul. 11, 1978; 4,192,012 to Schwee et al., entitled \"Crosstie Memory Bit Stretcher Detector,\" filed Nov. 8, 1978, patented Mar. 4, 1980; and 4,231,107 to Schwee et al., entitled \"Serriform Strip Crosstie Memory,\" filed Feb. 14, 1978, patented Oct. 28, 1980, all patents being assigned to the same assignee as the present invention.\nBriefly, Schwee et al., '770 disclose a polycrystalline thin film strip for storing digital information serially in the form of reversal domains. The reversal domains are propagated along the thin film strip, e.g., of permalloy, and then sensed to detect the stored digital information.\nSchwee, '659, discloses a serial access memory based on the propagation of crosstie walls and Bloch lines along domain walls in thin magnetic films. Variation of the current through conductors placed above the domain walls changes the fields along the walls causing the relocation of Bloch lines and crossties in the walls which causes propagation of the information contained in inverted Neel walls along the domain walls.\nSchwee et al., '107, discloses a crosstie memory using a plurality of permalloy thin film strips of uniform thickness each strip having parallel denticulated margins etched to align with the oblong axis thereof. Each pair of opposite denticles defines a distinct memory cell.\nThe other Patents listed above are cited to show additional prior art of crosstie memory development, and the development of certain types of detectors for use therewith. All of the listed patents are incorporated herein by reference.\nThe cited and discussed Patents do disclose certain versions and particular aspects of crosstie memories. But they do not disclose or teach, inter alia, the placing of magnetic domain walls into a matrix or array configuration using \"wiggle\" patterns thereby creating a random access memory element (memory array). Nor do the Patents teach a scheme to read and/or write in this memory array.\nThe prior art, as indicated hereinabove, teach some advances in crosstie memories including devices for detecting or reading-out binary information therefrom. However, insofar as can be determined, no prior art device or process incorporates the features and advantages of the present invention."}
{"text": "Embodiments of the present invention relate generally to aluminum alloy high pressure die castings and particularly to methods of enhancing mechanical properties of aluminum alloy high pressure die castings and to methods of manufacturing aluminum alloy high pressure die castings in high pressure die casting and heat treatment processes.\nHigh pressure die casting (HPDC) processes are widely used for mass production of metal components because of the processes' low cost and the close dimensional tolerances (near-net-shape) and smooth surface finishes they provide to the castings formed therefrom. For example, manufacturers in the car industry use HPDC to produce near-net-shape aluminum alloy castings for engine and, in particular, transmission applications.\nOne disadvantage of conventional HPDC processes, however, is that the HPDC castings generally are not amenable to solution treatment (T4) at high temperatures, such as about 500° C., for most high pressure die cast aluminum alloys. This significantly reduces the potential for precipitation hardening in the castings through a full T6 and/or T7 (=T4+T5, see detailed description below) heat treatment. The castings generally are not amenable to solution treatment (T4) due to a high quantity of porosity and voids in the components. The porosity and voids generally are attributable to shrinkage of the alloy from a low density liquid metal to a high density solid casting during solidification and, in particular, to gases, such as air, hydrogen or vapors formed from the decomposition of die wall lubricants, entrapped while filling the die with the molten metal. As such, virtually all HPDC castings have large gas bubbles formed therein. Further, internal pores containing gases or gas forming compounds within HPDC castings typically expand during conventional solution treatment at elevated temperatures, thereby, forming surface blisters on the castings. The presence of these blisters affects not only the appearance of castings, but the dimensional stability and, in particular, the mechanical properties of the castings as well.\nTherefore, to avoid the potential for blister formation, conventional aluminum alloy HPDC castings generally are used in as-cast and/or, to a lesser extent, in aged conditions, such as T5. Even with subjecting HPDC castings to a conventional T5 aging, however, the increase of yield strength, and other mechanical properties, is still very limited, since, in conventional as-cast aluminum alloy high pressure die castings, the concentrations of solutes available for strengthening in artificial aging (T5) are very low due to slow cooling after solidification. Additionally, the conventional single step isothermal aging (T5) at an intermediate temperature in many cases cannot maximize the mechanical properties for given concentrations of solutes in the material prior to aging. As a result, the mechanical properties of the conventional HPDC castings are usually low for a given composition of the aluminum alloy in comparison with other casting processes since the aluminum alloy castings made by other casting processes generally may be heat treated in full T6 or T7 conditions.\nDeveloped technologies, such as the use of vacuum to remove air in mold cavities during die filling, improve the quality of HPDC castings and their solution treat-ability. Use of these technologies, however, is still limited due to the high cost of facility and maintenance and operational complexity. Further, it also has been disclosed that blistering can be avoided, to a certain degree, by using much shorter solution treatment times and lower temperatures. For example, experiments with strengthening aluminum alloys 360 (Al-9.5Si-0.5Mg) and 380 (Al-8.5Si-3.5Cu) have shown that significant responses to aging are still possible following such modified solution treatments ([1] R. N. Lumley, R. G. O'Donnell, D. R. Gunasegaram, M. Givord, International Patent Application PCT/2005/001909; [2] R. N. Lumley, R. G. O'Donnell, D. R. Gunasegaram, M. Givord: Mat. Sci. Forum, 2006, vols. 519-522, pp. 351-359; [3] R. N. Lumley, R. G. O'Donnell, D. R. Gunasegaram, M. Givord: Proc 13th Die Casting Conference of the Australian Die Casting Association, Melbourne, Australia, 2006, P25; and [4] R. N. Lumley, R. G. O'Donnell, D. R. Gunasegaram, M. Givord, Metall Trans. 2008 in press). It appears, however, that these experiments are limited in value not only because the data disclosed merely are based on test specimens having very low porosity, but also because the solution (T4) heat treatment process parameter window is too narrow for highly complex HPDC castings.\nConventional T6 and/or T7 heat treatment processes for aluminum alloy castings normally involve following three stages: (1) solution treatment at a relatively high temperature below the melting point of the castings (also defined as T4), often for times exceeding 5 hours to dissolve its alloying (solute) elements and homogenize or modify the microstructure; (2) rapid cooling, or quenching, such as into cold or hot water, to retain the solute elements in a supersaturated solid solution; and (3) artificial aging (T5) by holding the casting for a period of time at an intermediate temperature suitable for achieving strengthening through precipitation. Solution treatment (T4) serves generally three main purposes: (1) dissolution of elements that lead to age hardening, (2) spherodization of un-dissolved particles and/or phases, and (3) homogenization of solute concentrations in the material. Quenching after T4 solution treatment retains the hardening solutes in a supersaturated solid solution (SSS) and creates a supersaturation of vacancies that enhances the diffusion and dispersion of precipitates. To maximize the yield strength, and other mechanical properties, of the casting, the precipitation of all strengthening phases should be prevented during quenching. Aging (T5, either natural or artificial) enables a controlled dispersion of strengthening precipitates. FIG. 1 shows a typical conventional T6 and/or T7 heat treatment cycle of an aluminum alloy.\nWith T5 aging (FIG. 1), there generally are three types of aging conditions, which are commonly referred as under-aging, peak-aging and over-aging. At an initial stage of aging, or pre-aging, Guinier-Preston (GP) zones and fine shearable precipitates form and the casting is considered to be under-aged. In this condition, mechanical properties of the casting usually are low. Increased time at a given temperature or aging at a higher temperature further evolves the precipitate structure and increases mechanical properties, such as yield strength, to a maximum levels to achieve the peak-aging/strength condition. Further aging decreases the mechanical properties and the casting becomes over-aged due to precipitate coarsening and its transformation of crystallographic incoherency. Simply for exemplary purposes, FIG. 2 shows an example of aging responses of cast aluminum alloys A356/357 manufactured under conventional sand casting processes and aged at a temperature of 170° C. For the period of aging time tested at a given aging temperature, the castings, whether sand castings or high pressure die castings referred to herein, undergo under-aged, peak-aged, and over-aged conditions.\nConsidering that conventional aluminum alloy HPDC castings generally inevitably contain internal porosity, artificial aging (T5) may be one of the ideal means (solutions) to help to achieve the desired mechanical properties in the castings without creating blisters. The strengthening resulting from aging occurs because the retained hardening solutes in the supersaturated solid solution form precipitates that are finely dispersed throughout the grains and that increase the ability of the casting to resist deformation by slip and plastic flow. Maximum strengthening may occur when the aging treatment leads to the formation of a critical dispersion of at least one type of these fine precipitates.\nIn addition, in conventional HPDC casting processes, the castings often are slowly cooled to a low temperature, such as below 200° C., prior to removal from the die to quench. This slow cooling to a low temperature significantly diminishes the subsequent aging potential of the casting since the hardening solute solubility decreases dramatically with the decrease in temperature, i.e., the lower the temperature, the lower the solubility. For example, the solubility of magnesium (Mg) in HPDC aluminum alloy A380 is about 0.34% at about 500° C. and decreases to nearly zero at about 200° C. Therefore, the conventional aluminum alloy high pressure die casting processes are ineffective in terms of energy consumption and achievable mechanical properties."}
{"text": "Many modern portable communications devices have a vehicle adapter that is mounted on a vehicle dashboard. The dashboard configuration provides access to power sources for battery charging and facilitates user access to device displays for operating the devices and interacting with applications running on the devices. Due to simultaneous charging, application usage (for example vehicle navigation), and/or prolonged exposure to direct sunlight, the portable communications devices or components therein may experience a significant rise in temperature.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of embodiments.\nThe apparatus and method components have been represented where appropriate by conventional symbols in the drawings, showing only those specific details that are pertinent to understanding the embodiments so as not to obscure the disclosure with details that will be readily apparent to those of ordinary skill in the art having the benefit of the description herein."}
{"text": "In the art of tear shades, flexible shade material (usually thin polyvinyl chloride) is wound onto a shade roller which is equipped in one of various possible ways to telescope within itself to shorten its length. The shade material is scored from top to bottom of the shade along one margin with spaced parallel score lines usually 1/2\" apart. When a width adjustment is desired, the user starts the material tearing along a selected score line, peels off the unwanted part, adjusts the length of the shade roller, and then mounts the adjusted shade in place.\nThe problem encountered in many arrangements is that the material tends not to start tearing accurately along a selected score line. It will be understood that the score lines must not be so deep that the shade material will tear easily otherwise there will be a risk that they may tear accidentally in use. For this reason it is preferable to have score lines which are on the shallow side, i.e. only sufficiently deep to provide uniform and accurate tearing. With such tear lines, however, the start of the tear becomes a problem.\nIn one form of tear shade a v-cut is formed along the bottom edge of the shade with the apex of the v on the tear line. This form has several drawbacks. First, starting the tear is so easy that there is a substantial risk of accidental tearing. Second, the apex of the v must be vey accurately positioned or else a tear will be started elsewhere than on the desired tear line.\nAccordingly, an object of the present invention is to provide convenient, sure means for starting the tear, in a tear shade along a selected score line."}
{"text": "1. Field of the Invention\nThis invention relates to an auto-threading method and apparatus suited to receive film and transfer the film to a tentering machine (crosswise stretching machine) in a biaxial stretching line.\n2. Description of the Prior Art\nIn general, the threading of band-like material such as film has hitherto been effected manually. This manual threading has to be done during a film-manufacturing operation, which leads to a highly dangerous situation. A velocity at which the film is manufactured has to be decreased to secure safety during the threading operation. This deteriorates the working efficiency. An additional drawback is that two or more workers are required. Immediately after the film is threaded, it cannot be connected but remains in a split state. In the meantime, a scrap process has to be performed at a stenter outlet. The split scrap films are non-uniformly discharged from the stenter outlet while being seized by the clips. Especially when manufacturing thick films, the films gather formlessly. It is required that the gathered films be released from the clips and then treated.\nAn apparatus disclosed to obviate the conventional defects described above is shown in FIGS. 10 to 12. Referring to these Figures, the numerals 03, 05 represent guide rolls; 04 a nip roll; 06, 07 draw-off rolls for drawing off the film; and 08, 09 guide plates, formed of a material exhibiting a small frictional coefficient, for guiding the film. An upper edge 010 of a cutter is attached to the front end of the guide plate 09 so that the upper part of the upper edge coincides with the upper face of the guide plate 09, or alternatively, the upper edge is separately fixed from the guide plate. Another guide plate 011 formed of the material exhibiting the small frictional coefficient is movable up and down with the aid of an air cylinder 013, wherein the rotational axis of the guide roll 03 serves as the center of movement for plate 011. A lower edge 012 of the cutter is adjustably secured to the top end of the other guide plate 011 so that the upper part of the lower edge coincides with the upper face of the other guide plate 011. A sprocket wheel 014 functions to drive a clip chain 018 and clips 02. A disk 015 pushes blades of the clips to bite the film. Reference numeral 016 denotes a shaft, and 017 a nut.\nThe film 01 is carried via the guide rolls 03, 05 to the draw-off rolls 06, 07 by some means, e.g., a carrier rope, wherein the film is drawn off. The air cylinder 013 is operated to move upwards toward the other guide plate 011. The film is cut off by the cutter fitted to the top end of the guide plate 011 in cooperation with the cutter of the guide plate 09. The film is held in such a position that the upper face of the other guide plate 011 coincides with the upper face of the guide plate 09, or the other guide plate 011 is positioned slightly higher than the plate 09.\nNext, the film 01 fed by the guide roll 03 and the nip roll 04 passes between the guide plates 08, 09. The film 01 reaches the clips 02 and is caught by these clips, and then the normal operation begins. Thereafter, the air cylinder 013 reverts to its original position preparatory to film threading.\nIn the prior art apparatus depicted in FIGS. 10 to 12, the film 01 passes between the guide plates 08, 09 with the help of the guide roll 03 and the nip roll 04, and both side edges at the end of the film are caught by the clips 02. However, when the soft or incompletely smooth film passes between the guide plates 08, 09, the end thereof remains free. Hence, if the film is pushed from behind, the film is bent, or creases are formed. It is therefore difficult for the clips to properly uniformly catch the film."}
{"text": "The invention relates to a retractable blade utility knife handle and more particularly to a retractable blade utility knife handle having a pair of releasably locking, pivotable handle halves.\nIn retractable-blade utility knives, the blade is slidably movable from a sheathed position to an extended unsheathed position wherein the knife blade projects through an opening in the knife handle to present a cutting edge. Such a utility knife is disclosed in Weimann U.S. Pat. No. 4,663,845, entitled \"Utility Knife\" wherein the knife incorporates a blade carrier which mounts and supports the blade within the interior of the knife handle for selective longitudinal movement therein. The blade carrier includes a thumb actuated button to release the blade carrier from one of several latching positions for slidably shifting the blade to another longitudinal position. The halves of the knife handle are connected by a central pivot. The handle halves are releasably locked so that they can be pivoted laterally relative to one another when the blade carrier is in any longitudinally shifted position, thereby allowing the blade to be replaced and providing access to a blade storage compartment.\nA typical shortcoming of retractable-blade utility knives is the lack of stability of the blade when the knife is forcibly twisted sideways during use. For example, under conditions in which a substantial force is applied to the side of the blade proximate the lower cutting edge, it is possible for the lower edge of the blade to slip off of the blade carrier, thereby freeing the blade to fall out of the handle. Another drawback of conventional utility knife handles is the possibility that the handles can be opened inadvertently. Yet another shortcoming of conventional utility knives is that the blade carrier and blade tend to rattle, as they are not fixed to the handle. A further drawback of conventional retractable blade knives is that if the blade carrier is shifted longitudinally when no blade is installed in the blade carrier, the corners of the blade carrier, on occasion, will become stuck at various locations in the handle, causing the blade carrier to shift or slide irregularly."}
{"text": "This invention relates to novel carboxycycloalkylamino derivatives. The carboxycycloalkylamino derivatives of the present invention are modulators of the sphingosine-1-phosphate (S1P) receptors and have a number of therapeutic applications, particularly in the treatment of hyperproliferative and autoimmune diseases, in mammals, especially humans, and to pharmaceutical compositions containing such compounds.\nThe S1P receptors 1-5 constitute a family of seven-transmembrane G-protein coupled receptors. These receptors, referred to as S1P1 to S1P5, are activated via binding by sphingosine-1-phosphate, which is produced by the sphingosine kinase phosphorylation of sphingosine. S1P receptors are cell surface receptors involved in a variety of cellular processes, including cell proliferation and differentiation, cell survival, and cell migration. S1P is found in plasma and a variety of other tissues and exerts autocrine and paracrine effects.\nRecent studies indicate that S1P binds to the S1P1 receptor to promote tumor angiogenesis by supporting the migration, proliferation and survival of endothelial cells (ECs) as they form new vessels within tumors (tumor angiogenesis) (Lee et al., Cell. 99:301-312 (1999) Paik et al., J. Biol. Chem. 276:11830-11837 (2001)). Because S1P is required for optimal activity of multiple proangiogenic factors, modulating S1P1 activation may affect angiogenesis, proliferation, and interfere with tumor neovascularization, vessel maintenance and vascular permeability.\nOther diseases or conditions that may be treated with the compounds of the present invention include organ transplant rejection and inflammatory diseases, which are believed to proceed via modulating the S1P receptors.\nThus, the identification of compounds which modulate the activity of the S1P1 receptor to regulate and modulate abnormal or inappropriate cell proliferation, differentiation, or metabolism is therefore desirable."}
{"text": "1. Field of the Invention\nThis invention relates to nozzles that disperse a fluid or fluids to a surface for cleaning, washing, blasting, or allied processes in which fluid impact with the surface is important. In addition, these nozzles have the ability to provide a spray over a large area with a liquid droplet size larger than conventional fantype nozzles having the same pressure and flow. Consequently, a spray over a large area with low overspray and atomization is obtained.\n2. Description of Related Art\nThere has been long-term interest in the use of pressurized fluids to impact surfaces. An example of one such application of pressurized fluids is the use of pressurized water for cleaning and washing of cars, trucks, industrial equipment, floors, driveways, and buildings. In any cleaning operation there are three functions to be performed: (1) the application of water or water and chemicals to soak dirt and film on the surface to be cleaned (soaking function), (2) the removal of dirt and film by the impact of the water jet (removal function), and (3) the application of water for rinsing the cleaned surface (rinsing function).\nFor a given amount of water in a given volume, the relative relationships between water pressure, velocity, flow rate, and impact energy are proportional: the higher the pressure, the higher the velocity; the higher the velocity, the higher the flow rate; the higher the velocity and flow rate, the higher the impact energy. However, the impact energy actually generated depends on the area of the surface impacted. This relationship between the impact energy and the area to be cleaned may be termed the impact energy density. To achieve a higher impact energy density, either the velocity and flow rate must be increased or the area impacted must be decreased.\nIn the soaking and rinsing functions, the flow rate of water, and thus the velocity and water pressure, must be sufficiently large to apply the necessary amount of water to cover the surface to be cleaned and to do so in a given amount of time. Particularly for the rinsing function, there is a minimum flow rate that is efficient in terms of time and water usage. In addition, the water pressure must be sufficiently large to project the water to the surface to be cleaned at a high velocity so that the impact energy of the water will be sufficient to dislodge dirt and other particles to perform the removal function. For water usage to be the most economical, a balance between the flow rate and the impact energy just be achieved. This balance must also be taken into account for each of the three cleaning functions.\nTo produce the desired effects during the soaking and rinsing functions, it is desirable to spread the water or water and chemicals over a large area. One common means of producing a flow over a large area is by a fan-type nozzle. The fan-type nozzle uses a small opening to limit the flow rate and expand the jet over a large area. The small opening causes the jet to break up into small droplets. The velocity of these droplets decreases as they impact the air. This decreased velocity means that the fan-type jet has a low impact energy. Furthermore, because the fan-type jet is spread over a large area, the impact energy density is low.\nTo produce the desired effects during the removal function, it is desirable to impact the water or water and chemicals over a small area so that there is a high impact energy density. One means of producing a flow over a small area is by using a nozzle that generates a zero-degree jet. Such nozzles are well-known to the art. A zero-degree jet is a jet that does not expand radially with respect to the direction of travel as it is projected from the nozzle. Because the droplets in a zero-degree jet follow the same path, the effects of air drag are decreased and the jet retains much more of its initial velocity than does a fan-type jet. Thus, the impact energy of a zero-degree jet is larger than that of a fan-type jet for two reasons. First, in contrast to a fan-type jet, a zero-degree jet impacts a smaller area, and thus, the impact energy density of a zerodegree jet is larger than that of a fan-type jet. Second, because the aerodynamic drag affects the fantype jet more, the fan-type jet loses its momentum more drastically as a function of distance travelled. Consequently, the zero-degree jet produces a larger impact energy and a larger impact energy density than a fan-type jet.\nAnother use of pressurized fluids is the application of chemicals such as insecticides and herbicides to a selected area. In these applications, it is important to direct the chemicals to the target area with a minimum of direct overspray or atomization of liquid to avoid susceptibility to drift. Consequently, in addition to requiring high pressure (for distance) and low flow, this application requires large droplet size which is inconsistent with that provided by conventional fan-type nozzle configurations."}
{"text": "1. Field of the Invention\nThis invention relates to an optical fiber amplifier and a dispersion compensating fiber module for use with an optical fiber amplifier.\n2. Description of the Related Art\nIn recent years, research and development of an optical communication system has been and is being performed energetically, and the importance of booster amplifiers, repeaters or preamplifiers which make use of the technique of optical amplification in which an erbium (Er) doped fiber (an erbium-doped-fiber may be hereinafter referred to as xe2x80x9cEDFxe2x80x9d) is employed has become apparent.\nFurther, due to the appearance of optical amplifiers, attention is drawn to an optical-amplifier-repeated transmission system since the transmission system plays a very important role in achievement of economization of a communication system in the multimedia society.\nBy the way, in an ordinary rare earth doped fiber optical amplifier which particularly amplifies a wavelength of a signal, the length of the doped fiber is set to a value at which a maximum gain is obtained in order to assure a high conversion efficiency from pump power to signal power.\nMeanwhile, in a wavelength division multiplexing (WDM) optical amplifier which amplifies many channels at the same time, it is important to keep the wavelength dependency of the gain as flat as possible. As a result, the rare earth doped fiber (which will be hereinafter discussed in connection with a representative EDF) must operate in a condition wherein the degree of the saturation of the gain is low. To this end, where the concentration of high level ions is represented by N2 while the concentration of all ions is represented by N1 and N2/N1 is defined as pump ratio, in order to raise the average pump ratio N2/N1 of the doped fiber over the entire length, the length of the doped fiber must be set short.\nHowever, if the doped fiber is formed short in this manner, then much residual pump power will leak out from the other end of the doped fiber, resulting in degradation of the conversion efficiency. Nevertheless, since the required pump power increases as the number of signal wavelengths increases, the output power of a semiconductor pump laser must be raised.\nIn particular, although it is apparent from the conservative law of energy that the pump power increases as the number of wavelengths increases, a wavelength multiplexing optical amplifier cannot be used in a condition in which it exhibits a high efficiency of conversion from pump power to signal power. This is because, since the rare earth doped fiber is intentionally formed short so as to prevent saturation in order to obtain a gain over a wide bandwidth or to make the gain flat, pump power which has not been converted into a signal will leak out from the other end of the doped fiber.\nAccordingly, while high pump power is required originally when comparing with ordinary amplification of only one signal channel, the rare earth doped fiber must be used in a condition wherein the pump power leaks out therefrom.\nThus, in order to effectively make use of thus leaking out residual pump light, a technique has been proposed wherein a reflecting mirror is provided at the other end of a doped fiber so that residual pump light is reflected by the reflecting mirror so as to be introduced back into the doped fiber so that it may be used for optical amplification again. The technique is disclosed in Japanese Patent Laid-Open Application No. Heisei 3-25985 or Japanese Patent Laid-Open Application No. 3-166782.\nHowever, where residual pump light is reflected by the reflecting mirror in this manner, the pump light is returned not only to the doped fiber but also to the pump source. This pump light may possibly give rise to unstable operation of the pump source such as interference.\nBy the way, while, due to the appearance of optical amplifiers, attention is drawn to an optical-amplifier-repeated transmission system which includes a plurality of repeating and amplifying optical amplifiers since it plays a very important role in achievement of economization of a communication system in the multimedia society as described above, the transmission system has subjects to be solved in terms of the dispersion compensation, reduction in nonlinear effects (effects having a bad influence on the transmission quality) in an optical fiber serving as a transmission line and economic wide bandwidth wavelength multiplexing transmission.\nGenerally, an optical fiber serving as a transmission line has a dispersion characteristic and accumulates a dispersion amount in proportion to the length thereof. Usually, however, in an optical fiber transmission system which employs regenerative repeaters, the dispersion amount is reset at the regenerative repeaters. Consequently, the accumulation of the dispersion amount does not make a problem.\nHowever, in an optical-amplifier-repeated transmission system, since a transmitted optical signal is repeated by a kind of analog amplification, the dispersion amount is accumulated. Accordingly, in order to eliminate the accumulation, the signal wavelength used for transmission should be set to a zero dispersion wavelength. This, however, provides the following subjects to be solved:\n1-1) Optical fibers have already been laid by a large amount, and unfortunately, those optical fibers have a zero dispersion wavelength at 1.3 xcexcm while an optical amplifier which is expected to be put into practical use soon can amplify only a signal of the 1.55 xcexcm band;\n1-2) It has been reported recently that, even if optical fibers whose zero dispersion wavelength is 1.55 xcexcm are laid newly to transmit a signal of 1.55 xcexcm, nonlinear effects occur actively in the optical fibers. This signifies that, if a signal wavelength equal to a zero dispersion wavelength is used for transmission, then undesirable nonlinear effects occur; and\n1-3) Particularly in wavelength multiplexing transmission, since a plurality of different signal wavelengths are involved, the concept that the signal wavelengths are set equal to a zero dispersion wavelength cannot be applied.\nAccordingly, it has been proposed recently to intentionally displace the signal wavelength from the zero dispersion wavelength suitably and compensate for the dispersion, for example, at the repeater.\nWhile research of dispersion compensators has been and is being performed actively in recent years in this manner, one of dispersion compensators which is expected to be most likely put into practical use is a dispersion compensating fiber (which may be referred to as xe2x80x9cDCFxe2x80x9d; here the term DCF is the abbreviation of Dispersion Compensating Fiber). The DCF, however, has the following subjects to be solved:\n2-1) Where fibers (transmission lines) laid already are utilized, a dispersion compensating fiber must be interposed as a device at each repeating point in order to perform dispersion compensation collectively at such each repeating point. Therefore, research and development is being directed to reduction in length of dispersion compensating fibers.\n2-2) When fibers are to be laid newly, it is a possible idea not to interpose a dispersion compensating fiber as a device but to lay a dispersion compensating fiber as part of a transmission line. For example, a transmission line of 40 km may be formed from a fiber of 20 km and a dispersion compensating fiber of 20 km. However, research and development of such a novel dispersion compensating fiber as just mentioned makes overlapping development with research and development of a dispersion compensating fiber for the application described in paragraph 2-1) above.\nIn summary, in wavelength multiplexing transmission, a wavelength dispersion must be compensated for, and since the compensation for a wavelength dispersion is expected to be most likely put into practical use where a dispersion compensating fiber is employed, it is prospective to use a dispersion compensating fiber. Further, it is investigated to incorporate a dispersion compensating fiber as a part into an optical amplifier repeater. Generally, however, the mode field diameter of a dispersion compensating fiber (DCF) is set small in order to compensate for a dispersion, and consequently, nonlinear effects are liable to occur and, as the dispersion amount to be compensated for increases, also the loss increases.\nThus, it is a possible method to compensate also for the loss of a dispersion compensating fiber using an optical amplifier. In this instance, the loss must be compensated for so that a transmission optical signal may not be influenced by nonlinear effects which degrade the quality of a signal such as self-phase modulation (SPM) and cross-phase modulation (XPM) occurring in the dispersion compensating fiber. Accordingly, the possible method has a problem in that designing of a level diagram is difficult. Further, while a flat and wide optical amplification bandwidth is required for an optical amplifier for WDM, also a rare earth doped fiber optical amplifier has a wavelength dependency of the gain. Accordingly, there is a subject to be solved in that it is difficult to realize a flat and wide amplification bandwidth.\nMeanwhile, a rare earth doped fiber optical amplifier having a high gain sometimes suffers from unnecessary oscillations which are produced when it performs optical amplification. If such unnecessary oscillations are produced, the rare earth doped fiber optical amplifier operates but unstably.\nFor example, in an erbium-doped-fiber optical amplifier, spontaneous emission light (ASE) of 1.53 to 1.57 xcexcm in wavelength is generated when optical amplification is performed, and since the ASE is repetitively reflected from reflection points in the erbium-doped-fiber optical amplifier, unnecessary oscillations are liable to be produced. Particularly with an erbium-doped-fiber optical amplifier adjusted for multiple wavelength collection amplification (that is, an erbium-doped-fiber optical amplifier having a high pump rate), since it has a high gain in the proximity of 1.53 xcexcm, unnecessary oscillations are liable to be produced at this wavelength. When such unnecessary oscillations are produced, the erbium-doped-fiber optical amplifier operates but unstably.\nIt is an object of the present invention to provide an optical fiber amplifier wherein stable operation of a pump source (pump light source) is assured and residual pump power which is produced when the average pump ratio is raised is utilized efficiently to improve the conversion efficiency.\nIt is another object of the present invention to provide an optical fiber amplifier and a dispersion compensating fiber module for an optical fiber amplifier employing a dispersion compensating fiber wherein the loss of the dispersion compensating fiber by Raman amplification can be compensated for making use of the fact that the threshold value of the Raman amplification is low because the mode field diameter of the dispersion compensating fiber is small.\nIt is a further object of the present invention to provide an optical fiber amplifier wherein, where a silica-type-optical-fiber having a Raman amplification function similarly to a dispersion compensating fiber is employed, the loss of the silica-type-optical-fiber by Raman amplification can be compensated for similarly to the case where a dispersion compensating fiber is used.\nIt Is a still further object of the present invention to provide an optical fiber amplifier which minimizes unstable operation of a rare earth doped fiber optical amplifier having a high gain or a rare earth doped fiber optical amplifier adjusted for multiple wavelength collective amplification.\nIn order to attain the objects of the present invention described above, according to an aspect of the present invention, there is provided an optical fiber amplifier including a rare earth doped fiber, which comprises first means for introducing pump light into one end of the rare earth doped fiber by way of a first optical coupler, second means for demultiplexing residual pump light originating from the pump light introduced into the one end of the rare earth doped fiber by the first means and arriving at the other end of the rare earth doped fiber by a second optical coupler and reflecting the demultiplexed residual pump light by reflection means so as to be introduced back into the rare earth doped fiber, and third means for preventing the residual pump light introduced back into the rare earth doped fiber by the second means from being introduced into a pump source, from which the pump light to be introduced into the rare earth doped fiber by the first means is produced, by optical isolation means so as to prevent unstable operation of the pump source.\nIn the optical fiber amplifier, when pump light is introduced into the one end of the rare earth doped fiber by way of the first optical coupler, residual pump light arrives at the other end of the rare earth doped fiber and is then demultiplexed by the second optical coupler, whereafter it is reflected by the reflection means so that it is introduced back into the rare earth doped fiber. In order to prevent unstable operation of the pump source caused by interference of the residual pump light introduced back into the rare earth doped fiber, the optical isolation means is interposed between the pump source and the first optical coupler. Consequently, the optical fiber amplifier is advantageous in that it makes use of the pump power with a high efficiency while assuring stabilized operation of the pump source.\nAccording to another aspect of the present invention, there is provide an optical fiber amplifier including a rare earth doped fiber, which comprises a pump source, a first optical coupler for introducing pump light from the pump source into one end of the rare earth doped fiber, a second optical coupler for demultiplexing residual pump light originating from the pump light introduced into the one end of the rare earth doped fiber by way of the first optical coupler and arriving at the other end of the rare earth doped fiber, a reflecting mirror for reflecting the residual pump light demultiplexed by the second optical coupler so as to be introduced back into the rare earth doped fiber by way of the second optical coupler, and an optical isolator interposed between the pump source and the first optical coupler for preventing unstable operation of the pump source arising from interference of the residual pump light introduced back into the rare earth doped fiber.\nIn the optical fiber amplifier, when pump light is introduced into the one end of the rare earth doped fiber by way of the first optical coupler, residual pump light arrives at the other end of the rare earth doped fiber and is demultiplexed by the second optical coupler, whereafter it is reflected by the reflecting mirror so that it is introduced back into the rare earth doped fiber. In order to prevent unstable operation of the pump source caused by interference of the residual pump light introduced back into the rare earth doped fiber, the optical isolator is interposed between the pump source and the first optical coupler. Consequently, the optical fiber amplifier is advantageous in that it makes use of the pump power with a high efficiency while assuring stabilized operation of the pump source.\nAccording to a further aspect of the present invention, there is provided an optical fiber amplifier including a first rare earth doped fiber and a second rare earth doped fiber disposed at front and rear stages, which comprises first means for introducing pump light into one end of one of the first rare earth doped fiber and the second rare earth doped fiber by way of an optical circulator having three or more ports and a first optical coupler, second means for demultiplexing residual pump light originating from the pump light introduced into the one end of the one rare earth doped fiber by the first means and arriving at the other end of the one rare earth doped fiber by a second optical coupler and reflecting the demultiplexed residual pump light by reflection means so as to be introduced back into the one rare earth doped fiber, and third means for causing the residual pump light reflected from the reflection means and introduced back into the one rare earth doped fiber by the second means to follow, after passing the one rare earth doped fiber, a different optical path by the optical circulator and multiplexing the residual pump light in the different optical path with an output of the other one of the first rare earth doped fiber and the second rare earth doped fiber by a third optical coupler.\nIn the optical fiber amplifier, pump light is first passed through the optical circulator having three or more ports and then introduced into the one end of the rare earth doped fiber at the front stage or the rear stage by the first optical coupler. Then, residual pump light originating from the pump light and arriving at the other end of the rare earth doped fiber is demultiplexed by the second optical coupler and then reflected by the reflection means so that it is introduced back into the rare earth doped fiber. The residual pump light is then introduced, after passing the rare earth doped fiber, into the different optical path by the optical circulator and is multiplexed with an output of the other rare earth doped fiber by the third optical coupler. Consequently, the optical fiber amplifier of the two stage construction just described is advantageous in that it makes use of the pump power with a high efficiency.\nAccording to a still further aspect of the present invention, there is provided an optical fiber amplifier including a first rare earth doped fiber and a second rare earth doped fiber disposed at front and rear stages, which comprises a pump source, a first optical coupler provided at one end of one of the first rare earth doped fiber and the second rare earth doped fiber, a second optical coupler provided at the other end of the one rare earth doped fiber, a third optical coupler provided at one end of the other one of the first rare earth doped fiber and the second rare earth doped fiber, a reflecting mirror for reflecting residual pump light demultiplexed by the second optical coupler so as to be introduced back into the one rare earth doped fiber by way of the second optical coupler, and an optical circulator having three or more ports connected to the pump source, the first optical coupler and the third optical coupler, and wherein pump light from the pump source is introduced into one end of the one rare earth doped fiber by way of the optical circulator and the first optical coupler, and residual pump light originating from the pump light introduced into the one end of the one rare earth doped fiber and arriving at the other end of the one rare earth doped fiber is demultiplexed by the second optical coupler and reflected by the reflecting mirror so as to be introduced back into the one rare earth doped fiber, whereafter the residual pump light is introduced, after passing the one rare earth doped fiber, into a different optical path by the optical circulator so that the residual pump light is thereafter multiplexed with an output of the other rare earth doped fiber by the third optical coupler.\nThe optical fiber amplifier of the two stage construction just described is advantageous in that it makes use of the pump power with a high efficiency.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier including a first rare earth doped fiber and a second rare earth doped fiber disposed at front and rear stages, which comprises first means for branching pump power at a ratio of n:1, n being a real number equal to or greater than 1, by an optical branching element, multiplexing the pump light from a port of the optical branching element by a first optical coupler and introducing the multiplexed light into one end of one of the first rare earth doped fiber and the second rare earth doped fiber, second means for extracting residual pump power originating from the pump light introduced into the one end of the one rare earth doped fiber by the first means and arriving at the other end of the one rare earth doped fiber by a second optical coupler connected to the other end of the one rare earth doped fiber, multiplexing the extracted residual pump power by a third optical coupler and introducing the multiplexed power into one end of the other one of the first rare earth doped fiber and the second rare earth doped fiber, and third means for multiplexing the pump power from another port of the optical branching element branched by the optical branching element and introducing the multiplexed power into the other end of the other rare earth doped fiber by a fourth optical coupler.\nIn the optical fiber amplifier, the pump power is branched at the ratio of n:1, and the pump light from a port of the optical branching element is multiplexed by the first optical coupler and then introduced into the rare earth doped fiber at the front stage or the rear stage. Then, residual pump power is extracted by the second optical coupler connected to the other end of the rare earth doped fiber and is then multiplexed by the third optical coupler. Then, the output light of the third optical coupler is introduced into the one end of the other rare earth doped fiber. Meanwhile, the branched pump power from another port of the optical branching element is introduced into the other end of and multiplexed in the other rare earth doped fiber by the fourth optical coupler. Consequently, the optical fiber amplifier of the two stage construction just described is advantageous in that it makes use of the pump power with a high efficiency.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier including a first rare earth doped fiber and a second rare earth doped fiber disposed at front and rear stages, which comprises a pump source, an optical branching element for branching pump power from the pump source at a ratio of n:1, n being a real number equal to or greater than 1, a first optical coupler for multiplexing the pump light from a port of the optical branching element and introducing the multiplexed light into one of the first rare earth doped fiber and second rare earth doped fiber, a second optical coupler for extracting residual pump power outputted from the one rare earth doped fiber, a third optical coupler for multiplexing the residual pump power extracted by the second optical coupler and introducing the multiplexed power into the other one of the first rare earth doped fiber and the second rare earth doped fiber, and a fourth optical coupler for multiplexing the pump power from another port of the optical branching element branched by the optical branching element and introducing the multiplexed power into the other rare earth doped fiber.\nIn the optical fiber amplifier, the pump power is branched at the ratio of n:1, and the pump light from a port of the optical branching element is multiplexed by the first optical coupler and then introduced into the rare earth doped fiber at the front stage or the rear stage. Then, residual pump power is extracted by the second optical coupler connected to the other end of the rare earth doped fiber and is then multiplexed by the third optical coupler. Then, the output light of the third optical coupler is introduced into the one end of the other rare earth doped fiber. Meanwhile, the branched pump power from another port of the optical branching element is introduced into the other end of and multiplexed in the other rare earth doped fiber by the fourth optical coupler. Consequently, the optical fiber amplifier of the two stage construction just described is advantageous in that it makes use of the pump power with a high efficiency.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier including a rare earth doped fiber, which comprises a pump source, an optical circulator having three or more ports one of which is connected to the pump source, a first optical coupler for multiplexing pump light introduced thereto from the pump source by way of the optical circulator and introducing the multiplexed light into one end of the rare earth doped fiber, a second optical coupler for demultiplexing residual pump light originating from the pump light introduced into the one end of the rare earth doped fiber by the first optical coupler and arriving at the other end of the rare earth doped fiber, a reflecting mirror for reflecting the residual pump light demultiplexed by the second optical coupler so as to be introduced back into the rare earth doped fiber by way of the second optical coupler, a residual pump light detector for detecting the residual pump light introduced back into the rare earth doped fiber by the reflecting mirror and inputted from the one end of the rare earth doped fiber to the optical circulator by way of the first optical coupler, and a controller for controlling the pump source so that the residual pump light detected by the residual pump light detector may be constant.\nIn the optical fiber amplifier, pump light is first passed through the optical circulator having three or more ports and is then introduced into the one end of the rare earth doped fiber by the first optical coupler. Then, residual pump light originating from the pump light and arriving at the other end of the rare earth doped fiber is demultiplexed by the second optical coupler and then reflected by the reflecting mirror so that it is introduced back into the rare earth doped fiber. The residual pump power comes out from the one end of the rare earth doped fiber and is then inputted by way of the first optical coupler to the optical circulator, by which it is introduced into the different optical path so that it is monitored by the residual pump light detector. Then, the residual pump power is kept constant under the control of the controller. Consequently, the wavelength characteristic of the gain of the optical fiber amplifier can be controlled so that it may not be varied irrespective of any variation of the input level. Consequently, the optical fiber amplifier is advantageous in that it can be realized readily as a multiple wavelength collective amplifier.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber optical amplification element formed from a rare earth doped fiber, and a Raman optical amplification element which is pumped with pump light to cause Raman amplification to occur, the rare earth doped fiber optical amplification element and the Raman optical amplification element being connected in cascade connection.\nThe optical fiber amplifier is advantageous in that it makes use of the pump power with a high efficiency while it has a two stage construction.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber optical amplification element formed from a rare earth doped fiber, a Raman optical amplification element which is pumped with pump light, which is capable of pumping the rare earth doped fiber optical amplification element, to cause Raman amplification to occur, the rare earth doped fiber optical amplification element and the Raman optical amplification element being connected in cascade connection, and a pump source for supplying pump light for pumping the rare earth doped fiber optical amplification element and the Raman optical amplification element.\nIn the optical fiber amplifier, since it includes the pump source for supplying pump light for pumping the rare earth doped fiber optical amplification element and the Raman optical amplification element, the pump power can be utilized with a high efficiency and the number of pump sources to be used can be reduced, which contributes to simplification in construction and reduction in cost.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber optical amplification element formed from a rare earth doped fiber, and a Raman optical amplification element formed from a dispersion compensating fiber which is pumped with pump light to cause Raman amplification to occur, the rare earth doped fiber optical amplification element and the Raman optical amplification element being connected in cascade connection at two front and rear stages.\nThe optical fiber amplifier just described is advantageous in that it makes use of the pump power with a high efficiency while it has a two stage construction.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber and a dispersion compensating fiber disposed at two front and rear stages, a first pump source for producing pump light of a first wavelength band for the rare earth doped fiber, a first optical coupler for introducing the pump light from the first pump source into the rare earth doped fiber, a second pump source for producing pump light of a second wavelength band for the dispersion compensating fiber, and a second optical coupler for introducing the pump light from the second pump source into the dispersion compensating fiber, the dispersion compensating fiber being pumped with the pump light of the second wavelength band from the second pump source to cause Raman amplification to occur.\nWith the optical fiber amplifier, compensation for the loss of the dispersion compensating fiber by Raman amplification can be achieved while optical amplification is performed by the rare earth doped fiber.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises an erbium-doped-fiber and a dispersion compensating fiber disposed at two front and rear stages, a pump source for producing pump light, and an optical coupler for introducing the pump light from the pump source into the erbium-doped-fiber, the dispersion compensating fiber being pumped with residual pump light from the erbium-doped-fiber to cause Raman amplification to occur.\nWith the optical fiber amplifier, compensation for the loss of the dispersion compensating fiber by Raman amplification can be achieved while optical amplification is performed by the erbium-doped-fiber.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises an erbium-doped-fiber and a dispersion compensating fiber disposed at two front and rear stages, a pump source for producing pump light, and an optical coupler for introducing the pump light from the pump source into the dispersion compensating fiber, the erbium-doped-fiber being pumped with residual pump light from the dispersion compensating fiber.\nWith the optical fiber amplifier, compensation for the loss of the dispersion compensating fiber by Raman amplification can be achieved while optical amplification is performed by the erbium-doped-fiber.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a dispersion compensating fiber doped with a rare earth element, a pump source for producing pump light for the dispersion compensating fiber, and an optical coupler for introducing the pump light from the pump source into the dispersion compensating fiber.\nWith the optical fiber amplifier, since the dispersion compensating fiber used is doped with a rare earth element, dispersion compensation can be performed by the dispersion compensating fiber, and the loss of the dispersion compensating fiber can be reduced simultaneously. The optical fiber amplifier with the dispersion compensation function is advantageous also in that it can optically amplify signal light sufficiently.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises an erbium-doped-fiber and a dispersion compensating fiber disposed at two front and rear stages, a pump source for producing pump light for the erbium-doped-fiber, an optical coupler for introducing the pump light from the pump source into the erbium-doped-fiber, and an optical filter interposed between the erbium-doped-fiber and the dispersion compensating fiber for intercepting residual pump light coming out from the erbium-doped-fiber.\nWith the optical fiber amplifier, leakage pump power Raman amplifies the dispersion compensating fiber. Consequently, the optical fiber amplifier is prevented from unstable operation or from variation of the wavelength dependency of the amplification band thereof.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber optical amplification element formed from a rare earth doped fiber, and a Raman optical amplification element formed from a silica-type-optical-fiber which causes, when pumped with pump light, Raman amplification to occur, the rare earth doped fiber optical amplification element and the Raman optical amplification element being connected in cascade connection at two front and rear stages.\nThe optical fiber amplifier just described is advantageous in that it makes use of the pump power with a high efficiency while it has a two stage construction.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a silica-type-optical-fiber and an erbium-doped-fiber provided at a front stage and a rear stage, respectively, a silica-type-optical-fiber pump source for producing pump light of a wavelength band for the silica-type-optical-fiber, an optical coupler for introducing the pump light from the silica-type-optical-fiber pump source into the silica-type-optical-fiber, an erbium-doped-fiber pump source for producing pump light of a wavelength band for the erbium-doped-fiber, and another optical coupler for introducing the pump light from the erbium-doped-fiber pump source into the erbium-doped-fiber, the silica-type-optical-fiber being pumped with the pump light from the silica-type-optical-fiber pump source to cause Raman amplification to occur.\nWith the optical fiber amplifier, compensation for the loss of the silica-type-optical-fiber by Raman amplification can be performed while optical amplification by the erbium-doped-fiber is performed.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises an erbium-doped-fiber having a low noise figure and a silica-type-optical-fiber provided at a front stage and a rear stage, respectively, a silica-type-optical-fiber pump source for producing pump light of a wavelength band for the silica-type-optical-fiber, an optical coupler for introducing the pump light from the silica-type-optical-fiber pump source into the silica-type-optical-fiber, an erbium-doped-fiber pump source for producing pump light of a wavelength band for the erbium-doped-fiber, and another optical coupler for introducing the pump light from the erbium-doped-fiber pump source into the erbium-doped-fiber, the silica-type-optical-fiber being pumped with the pump light from the silica-type-optical-fiber pump source to cause Raman amplification to occur.\nWith the optical fiber amplifier, compensation for the loss of the silica-type-optical-fiber by Raman amplification can be performed while optical amplification by the erbium-doped-fiber is performed.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber optical amplification element formed from a rare earth doped fiber and having a low noise figure, the rare earth doped fiber optical amplification element being disposed as a front stage amplification element, a Raman optical amplification element for causing Raman amplification to occur when pumped with pump light, the Raman optical amplification section being disposed as a middle stage amplification element, and another rare earth doped fiber optical amplification element formed from a rare earth doped fiber and disposed as a rear stage amplification element.\nWith the optical fiber amplifier, the compensation effect of the Raman optical amplification element can be increased. Consequently, a wide bandwidth optical amplifier can be realized while achieving simplification in structure and reduction in cost.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a first erbium-doped-fiber having a low noise figure, a dispersion compensating fiber and a second erbium-doped-fiber provided at a front stage, a middle stage and a rear stage, respectively, a first erbium-doped-fiber pump source for producing pump light of a wavelength band for the first erbium-doped-fiber, an optical coupler for introducing the pump light from the first erbium-doped-fiber pump source into the first erbium-doped-fiber, a dispersion compensating fiber pump source for producing pump light of a wavelength band for the dispersion compensating fiber, another optical coupler for introducing the pump light from the dispersion compensating fiber pump source into the dispersion compensating fiber, a second erbium-doped-fiber pump source for producing pump light of a wavelength band for the second erbium-doped-fiber, and a further optical coupler for introducing the pump light from the second erbium-doped-fiber pump source into the second erbium-doped-fiber, the dispersion compensating fiber being pumped with the pump light from the dispersion compensating fiber pump source to cause Raman amplification to occur.\nWith the optical fiber amplifier, the compensation effect of the Raman optical amplification element can be increased. Consequently, a wide bandwidth optical amplifier can be realized while achieving simplification in structure and reduction in cost.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a first erbium-doped-fiber having a low noise figure, a silica-type-optical-fiber and a second erbium-doped-fiber provided at a front stage, a middle stage and a rear stage, respectively, a first erbium-doped-fiber pump source for producing pump light of a wavelength band for the first erbium-doped-fiber, an optical coupler for introducing the pump light from the first erbium-doped-fiber pump source into the first erbium-doped-fiber, a silica-type-optical-fiber pump source for producing pump light of a wavelength band for the silica-type-optical-fiber, another optical coupler for introducing the pump light from the silica-type-optical-fiber pump source into the silica-type-optical-fiber, a second erbium-doped-fiber pump source for producing pump light of a wavelength band for the second erbium-doped-fiber, and a further optical coupler for introducing the pump light from the second erbium-doped-fiber pump source into the second erbium-doped-fiber, the silica-type-optical-fiber being pumped with the pump light from the silica-type-optical-fiber pump source to cause Raman amplification to occur.\nWith the optical fiber amplifier, the compensation effect of the silica-type-optical-fiber can be increased. Consequently, a wide bandwidth optical amplifier can be realized while achieving simplification in structure and reduction in cost.\nAccording to a yet further aspect of the present invention, there is provided a dispersion compensating fiber module for an optical fiber amplifier, which comprises a dispersion compensating fiber, and a pump source for pumping the dispersion compensating fiber to cause Raman amplification to occur.\nWhere an optical fiber amplifier is constructed using the module wherein the dispersion compensating fiber is pumped to cause Raman amplification to occur, it exhibits a reduced loss due to reduction of the loss by the dispersion compensating fiber.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier including a dispersion compensating fiber, which comprises a pump source, and an optical coupler for introducing pump light from the pump source into the dispersion compensating fiber, the dispersion compensating fiber being pumped with pump light from the pump source to cause Raman amplification to occur.\nAlso the optical fiber amplifier is advantageous in that the loss of the dispersion compensating fiber can be reduced.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a pump source, and an optical coupler for introducing pump light from the pump source into the silica-type-optical-fiber, the silica-type-optical-fiber being pumped with the pump light from the pump source to cause Raman amplification to occur.\nThe optical fiber amplifier is advantageous in that the loss of the silica-type-optical-fiber can be reduced.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises a rare earth doped fiber optical amplification element formed from a rare earth doped fiber, and an optical fiber attenuation element formed from an optical fiber or an optical fiber with an optical isolator for suppressing unstable operation of the rare earth doped fiber optical amplification element.\nThe optical fiber amplifier is advantageous in that stabilized optical amplification can be achieved with unstable operation of the rare earth doped fiber optical amplification element suppressed.\nAccording to a yet further aspect of the present invention, there is provided an optical fiber amplifier, which comprises an optical amplification unit including a front stage optical amplification element and a rear stage optical amplification element each formed as a rare earth doped fiber optical amplification element formed from a rare earth doped fiber, and an optical fiber attenuation element formed from an optical fiber or an optical fiber with an optical isolator interposed between the front stage optical amplification element and the rear stage optical amplification element of the optical amplification unit for suppressing unstable operation of the optical amplification unit.\nThe optical fiber amplifier is advantageous in that stabilized optical amplification can be achieved with unstable operation of the optical amplification unit suppressed.\nFurther objects, features and advantages of the present invention will become apparent from the following detailed description when read in conjunction with the accompanying drawings in which like parts or elements are denoted by like reference characters."}
{"text": "The present invention relates to memory cards, and relates more particularly to a contact terminal for memory card plug connectors which prevents contact errors.\nVarious memory cards have been developed for connection to a card jack on a computer system for storing data. As illustrated in FIG. 1, the card jack, referenced by 1, comprises an insertion slot 11 at the front side, and a set of contact pins 12 for connection to the circuit board of the computer system. The memory card, referenced by 3, has a plug connector 4 for insertion into the insertion slot 11, and a set of contact terminals 41 mounted in the plug connector 4 for connection to the contact pins 12 respectively. As illustrated in FIGS. 2 and 3, each of the contact terminals 41 comprises a curved upper clamping section 411 and a curved lower clamping section 412 for contacting one contact pin 12. Because the contact pin 12 is disposed in contact with the upper clamping section 411 and the lower clamping section 412 at a respective contact point, the contact pin 12 tends to displace, causing a contact error. FIGS. 4 and 5 shows another structure of contact terminal for this purpose. This contact terminal, referenced by 42, comprises two substantially C-shaped clamping portions 420 at two opposite ends, and three longitudinal connecting portions 421, 422, and 423 connected between the C-shaped clamping portions 420. However, this structure of contact terminal is still not satisfactory in function. During the installation of the memory card, much effort must be given to the plug connector of the memory card against the card jack so that the contact terminals 42 of the plug connector of the memory card can be respectively forced into engagement with the contact pins 12 of the card jack. However, when the contact terminals 42 are forced into engagement with the contact pins 12 of the card jack, the contact pins 12 tend to be deformed or broken. Another drawback of this structure of contact terminal is that frequently mounting and dismounting the contact terminal will cause the longitudinal connecting portions to lose their elastic resilient material property. If the longitudinal connecting portions lose their elastic resilient material property, a positive contact between the contact terminal and the respective contact pin cannot be obtained."}
{"text": "This application claims priority to German Application DE 101 48 696.0 filed on Oct. 2, 2001 and DE 102 23 296.2 filed on May 24, 2002.\n1. Technical Field of the Invention\nThe invention relates generally to self-propelled machines and, more particularly, to a method and apparatus for prolonging the life of a clutch on such machines, especially agricultural or utility machine.\n2. Description of the Related Art\nEP 0 728 962 and its counterpart, U.S. Pat. No. 5,695,422, describe, inter alia, a hydrostatic-mechanical torque-divided transmission for agricultural machines, and other utility machines. To change direction of the machine the rotation of the input shaft is reduced to zero. Typically this is accomplished, at least in part, by reducing fuel feed and allowing the engine to gradually slow the rotation of the input shaft. Torque is transferred via a clutch inside the transmission. In the case of high torque transfer, this leads to both considerable heat generation in the clutch and to premature wear. To reduce heat generation and wear-related drawbacks, it is known that several clutches or braking elements can be integrated in the respective drive train. But this still leads to elevated heat generation in the drive train and has the disadvantage that a considerable number of additional assemblies are required.\nThe present invention is directed to overcoming one or more of the problems set forth above.\nAn aspect of the invention is to provide a method and an apparatus for the control of a self-propelled machine which, with little structural expenditure, minimizes wear of the drive train, reduces heat generation in the drive train, reduces the high power idle that occurs with ordinary methods for changing the direction of the machine, and/or lowers overall fuel consumption of the self-propelled machine.\nAnother aspect of the present invention is to control a self-propelled machine which includes a gearbox with an output shaft, a clutch, and a braking system. When a change of speed and/or direction of rotation of the output shaft is required, a management controller effects the change by at least partially bypassing the clutch and engaging the machine braking system.\nWhen changes of speed and/or direction of rotation of the gearbox output shaft are assisted by the machine braking system, the braking function of the clutch is considerably reduced or completely excluded. This leads to a considerable reduction in wear of the clutch, so that its service life can be greatly increased, compared with known designs in the art. Also, apart from reducing the load on the clutch, a considerable reduction of engine idle power generated during deceleration is accomplished, so that the machine needs considerably less fuel, particularly in the case of changes in the direction of travel.\nTo not impair driving stability, it is of further advantage if associated with each land wheel or each axle of the machine are wheel brakes which are operated during changes of speed and/or direction of rotation of the gearbox output shaft.\nA particularly advantageous embodiment arises if the self-propelled machine includes a traction machine and towed equipment. In this embodiment, the land wheels and/or axles of both the traction machine and the towed equipment are wheel brakes. For changes of speed and/or direction of rotation of the output shaft, the wheel brakes of both the traction machine and the towed equipment are operated. In this way, even with tractor/trailer combinations, it is ensured that, in addition to the wear-minimization of the clutch of the traction machine, high driving stability of the combination is accomplished.\nTo obtain high flexibility of the machine braking system in the case of changes of speed and/or direction, an advantageous embodiment includes a machine brake management system.\nTo obtain high driving stability independently of the driver\"\"s wish to decelerate, machine-generated input signals delivered to the management controller can be used to monitor the driving stability and achieve the inputted wish to decelerate.\nIn accordance with the present invention there is provided an apparatus for prolonging the life of a clutch of a self-propelled machine having a drive train including an engine, a gearbox with an output shaft, and a machine braking system; means for generating an input signal relating to a function of the drive train; a management controller for receiving the input signal, analyzing the input signal for its effect on the drive train, and generating a plurality of output signals; and means for transmitting one of the output signals to the machine braking system and another of the output signals to the clutch, and for automatically and cooperatively engaging the machine braking system and disengaging the clutch to assist a change in the speed or direction of rotation of the output shaft.\nOther aspects, objects and advantages of the present invention can be obtained from a study of the drawings, the disclosure and the accompanying claims."}
{"text": "The present invention concerns a device for indexable coding for an electrical connector of the type comprising a correction component having a rear part comprising a right prism of which at least one of the lateral faces constitutes an angular positioning reference surface for indexing as well as an elongated front part extending in a longitudinal direction and having a correction profile.\nIndexable coding devices of this type are currently used in the connection field, in particular for connectors designed for aviation.\nIt is desirable to increase the number of correction positions, which is currently feasible only at the price of a certain complication, thus resulting in a relatively high cost."}
{"text": "1. Field of the Invention\nThe present invention relates to card edge connectors, more particularly to card edge connectors for mating with a memory card.\n2. Description of Related Art\nCard edge connectors are widely employed in computer device to receive a memory card, such as serial DDR3 and serial DDR4. The connector usually has an elongated housing, a plurality of contacts retained in the housing for electrically connecting a corresponding memory card, and two ejectors at both ends thereof for locking the memory card. The housing has a pair of sidewalls and a central slot for receiving the memory card. Each side wall defines a plurality of passageways extending therethrough. The passageways communicate with the center slot. The contacts are arranged in two rows retained in the passageways of two side walls respectively, and the contacts extend into the center slot to sandwich the memory card and ensure the electronic contact with the card. The memory card has a pair of notches at a middle position of two side edges thereof. The ejectors has a lock portion at an upper side thereof for locking the notches respectively to keep a stable working state, and an ejecting portion at a lower side thereof for ejecting the memory card to a convenient replacement.\nWith the development of the electronic industry, the number of the gold fingers increases to broaden the signal transmission channel and improve the signal transmission speed. Of course, the number of the contacts in the card edge connector increase codependent. When the memory card is inserted into the connector, it need overcome a big insertion force to insert the memory card into the central slot, therefore, we could not be aware of whether the card is in place or not, and it always makes inconvenience and confuse for users.\nHence, an improved card edge connector is desired to overcome the above problems."}
{"text": "2-benzylaniline is an important compound as an intermediate material for pharmaceuticals. One of the known methods of manufacturing this compound is to reduce 2-aminobenzophenone with metallic sodium or hydrazine (J. Chem. Soc., 292, 1948; Chem. Ber., 96, 765, 1963). However, the starting 2-aminobenzophenone used therefor is difficult to obtain, and further, the cost of 2-aminobenzophenone is relatively high. Therefore, there has been a demand to develop a less expensive manufacturing method for industrial use.\nAnother method is known which applies Friedel-Crafts reaction-to 2-aminobenzylchloride hydrochloride and benzene by using aluminum chloride (Chem. Ber., 61, 2276, 1928). Another known method is to reduce o-nitrodiphenyl methane obtained by Friedel-Crafts reaction using 2-nitrobenzylchloride and benzene with aluminum chloride (J. Am. Chem. Soc., 53, 1428, 1931). However, benzene is a carcinogenic material and its safety problem remains unsolved in industrial production."}
{"text": "1. Field of the Invention\nThis invention relates to a passenger passive restraint air bag module, and more particularly, to an improvement therein designed to minimize deformation thereof and also possible damage to the dashboard or instrument panel of an automotive vehicle during deployment of the bag. The improvement is in the structure internal to the module for housing and positioning the gas generator or inflator and the air bag and for directing the inflating gas flow or jet for best deployment of the bag.\n2. Description of the Prior Art\nSome terminology herein is used for convenience in reference only and is not intended to be limiting. The words \"forwardly\" and \"rearwardly\" refer to the normal forward and reverse directions of travel of the vehicle to which a passenger passive restraint module is attached. The word \"thrust neutral\" refers to the production by an inflator of zero thrust when accidentally initiated as, for example, during shipping, storage or handling. That is to say, the gas discharge openings in the inflator are so positioned that the gas is discharged in opposing directions whereby there are no resulting forces tending to cause physical movement of the inflator. Hence, the inflator will expend the energy generated thereby, generally in place.\nThe state of the prior art is indicated by the following U.S. patents:\n______________________________________ U.S. Pat. No. Issue Date Patentee ______________________________________ 3,414,292 December 3, 1968 S. Oldberg et al. 3,715,131 February 6, 1973 E. K. Hurley et al. 3,880,447 April 29, 1975 W. F. Thorn et al. 4,332,398 June 1, 1982 G. R. Smith 4,817,828 April 4, 1989 G. W. Goetz 4,832,300 June 27, 1989 J. F. Zionmek et al. 4,941,678 July 17, 1990 D. R. Lauritzen et al. ______________________________________\nThe Oldberg et al. patent discloses a safety device for providing protection for vehicle passengers comprising a folded inflatable crash bag closely surrounding an elongated cylindrical diffuser member, being secured thereto by clamps. The diffuser member defines a chamber in which an inflator comprising a fluid reservoir in an inner container is concentrically located. For providing a tight seal the inner container is seam welded at the opposite ends thereof to the diffuser member. When activated by explosive means, fluid from the reservoir flows rearwardly from one side only of the inflator against an adjacent inner wall of the diffuser member. The flow divides and is discharged forwardly through a single row of longitudinally spaced apart openings, that are provided on one side only of the diffuser member, against a fixed baffle member. The baffle member divides and redirects the fluid rearwardly into the inflatable bag.\nIn the fabrication and assembly of inflatable devices for protecting vehicle passengers, the inflator poses the greatest potential hazard. Accidental initiation thereof during shipping, storing and installation in the safety device could result in injury to personnel and also cause property damage, particularly when the inflator is not of the thrust neutral type. Thus, in order to minimize the possibility of such injury and damage, it is desirable to reduce the amount of handling to which the inflator is subjected during the fabrication and assembly of the safety device. Installation of the inflator as a last operation in the assembly of the device would provide a substantial reduction in the risk of injury and damage. With the use of an inflator of the thrust neutral type, the risk of such injury and damage would be further reduced.\nThe construction of the Oldberg et al. safety device is such as to preclude installation of the inflator as a last operation of the assembly of the device. Specifically, the inflatable bag is provided in closely surrounding relation to the diffuser member, with the inflator, that is, the inner container, welded at the ends thereof to the diffuser member. Moreover, with fluid flowing from the reservoir of fluid from one side only of the inflator and from one side only of the diffuser member, neither the inflator alone, nor the inflator and diffuser member, as a unit, are thrust neutral.\nThe Hurley et al. patent discloses a folded inflatable bag surrounding a gas generator having a head assembly of cylindrical shape provided with a closed end and an open aft end that is closed by a closure plate. A centrally positioned orifice that is normally closed by a rupture disc is provided in the closure plate. Enclosing the cylindrical portion of the head assembly is a concentrically positioned diffuser having a closed end in the shape of a dished head, thereby providing an annular cavity about the gas generator. Upon initiation of the gas generator, the rupture disc releases, in response to pressure in the head assembly exceeding the burst pressure of the disc, thereby allowing the generated gas to flow out of the orifice in the closure plate of the aft end of the gas generator and to impinge on the inner wall at the end of the diffuser. This causes the gas to be dispersed into the annular cavity and out of a plurality of openings spaced around the wall of the diffuser. With the generated gas flowing out of the aft end of the gas generator, the operation thereof is not thrust neutral. Nor does the Hurley construction allow fabriciation of the device with the installation of the gas generator as a last operation of the assembly.\nThe Thorn et al. patent discloses a folded inflatable bag positioned in close contact with an elongated cylindrical gas generator. The gas generator is provided with rows of gas outlet ports that extend 360.degree. around the cylinder. The gas is discharged in opposing directions and produces counteracting thrusts, and thus, is thrust neutral. A gas redirecting plate in the shape of a half cylindrical container is positioned in spaced relation around the forwardly facing portion of the gas generator. The plate redirects forwardly discharged inflating gas in a rearward direction into the inflatable bag. Redirected gas together with gas that directly is discharged into the inflatable bag deploy the bag, which normally is folded, to an inflated condition. In the stored position thereof, the bag is positioned closely adjacent to and in contact with the rearward half side of the gas generator. This precludes installation of the gas generator as a last operation in the assembly of the apparatus.\nThe Smith patent discloses an inflatable restraint system including a folded inflatable cushion, a concentric elongated cylinder inflator-diffuser combination, and an elongated horizontally orientated nozzle. The inflator and diffuser are located within the confines of the folded cushion and are connected as a unit to the side walls of a housing therefor by means of bolts which extend into tapped bosses, which bosses are rigid with the combined unit. Consequently, neither the combined unit, nor the inflator can be installed as a last operation in the assembly of the system.\nIn the Goetz patent an elongated cylindrical inflator and a folded air bag are mounted in a rigid reaction canister which is fixed to the dashboard of a vehicle, with the orientation and arrangement of gas discharge openings or exhaust ports in the inflator being such that, when activated, an initial flow of generated gas is directed rearwardly by the inflator to inflate the bag. Excess gas flows from the inflator in a forward direction into the surrounding environment. This result is achieved by using rupturable foil layers to close both forwardly and rearwardly directed inflator exhaust ports, with rupturing of the rearwardly directed ports occurring at a pre-selected pressure that is lower than that at which the forwardly directed ports rupture. Thus, the inflator is not thrust neutral, nor is it capable of installation as a last operation of the assembly of the apparatus because of the positioning in the reaction canister of the folded air bag in close proximity to the inflator.\nAs disclosed in the Goetz patent and the other aforementioned prior art patents, it is known in prior art installations to provide a container to house the air bag and the inflator. The container is generally referred to as a reaction canister and provides a medium for installing and retaining a module in the instrument panel by the utilization of suitable brackets.\nThe reaction canister provides protection for the inflator, the diffuser, if provided, and the air bag until the time of deployment of the bag. Additionally, the reaction canister absorbs the loads that are produced upon deployment of the bag, which loads, typically, are large. Unless sufficiently absorbed, they can cause serious damage to the interior of the vehicle, and in particular, to the instrument panel.\nA drawback to the use of a reaction canister in an air bag module for passenger protection is the envelope in the dashboard or instrument panel that is allotted for the installation of the module. Resulting restrictions, usually in height, causes lengthened, that is, deeper, top and bottom panels in order to accommodate the folded volume of the air bag. The bag is then restricted for unfolding and must travel rearwardly a distance greater than desirable before unfolding. Because of the large amount of gas that is produced and the distance the bag needs to travel before unfolding, pressure builds up in the reaction canister to a level that tends to expand the canister and cause it to bulge. This condition, known as bell mouthing, is very destructive of the instrument panel, and is particularly objectionable in low speed crashes where other damage to the vehicle is small.\nThe Zionmek et al. patent and the Lauritzen et al. patent disclose the use of a tether strap to resist the spreading forces on the reaction canister and thereby preclude bell mouthing upon deployment of the air bag. This technique allows the use of lighter weight and less expensive reaction canisters. The structural arrangement, however, leaves something to be desired in respect of the introduction of undesirable complication into the manufacturing and assembling operations, and moreover, does not allow installation of the inflator as a last operation in the assembly of the module.\nThus, there is a need and a demand for an improved reaction canister for use in an air bag module for the passenger side of a vehicle."}
{"text": "1. Technical Field\nThe present disclosure relates to camera systems, and particularly to an image stabilization system.\n2. Description of Related Art\nIn normal use of a camera, light rays from an object transmit into the camera and fall on a particular region of an image sensor. The image sensor forms an optical image associated with the object at a first position.\nHowever, camera-shake due to hand-held shooting during optical image capture can result in the camera moving slightly relative to the object. Consequently, light rays from the object may fall on a different region of the image sensor causing the optical image associated with the object at a second position to form a blurry optical image. Current image stabilization systems that deal with this problem are expensive and complicated.\nTherefore, a new image stabilization system is desired to overcome the above-mentioned problems."}
{"text": "The present disclosure relates generally to the field of network traffic data management and, more particularly, to a system and method for managing data prioritization and connections by a client.\nData networks operate by transferring data between various devices. The amount of data that can be handled simultaneously by a network is generally constrained by the network's bandwidth. Accordingly, high bandwidth applications, such as voice communications, may use much or all of a network's bandwidth. Furthermore, even if the network has unused bandwidth capacity in one area, a single link (e.g., to a user's computer) may be burdened due to heavy traffic on that link. This traffic burden becomes even more pronounced if multiple data streams (e.g., multiple conversations) are occurring simultaneously. Not only does this require network resources, but it also places additional stress on the user's computer.\nAccordingly, what is needed are an improved system and method for managing data traffic on a network."}
{"text": "1. Field of the Invention\nThe present invention relates to a reference voltage generating method, and in particular relates to a reference voltage device.\n2. Description of the Related Art\nIntegrated circuits often require a reference voltage which remains stable under PVT (process, voltage, temperature) variations. FIG. 1 depicts a conventional generating circuit of a bandgap voltage serving as a reference voltage for integrated circuits. Referring to FIG. 1, the generating circuit 1 comprises an operational amplifier OPAMP, bipolar devices B11 and B12, and resistors R11 to R13. The values of the resistors R11 and R12 are the same. The operational amplifier OPAMP of a high gain is used, and base-emitter voltage Vbe1 and Vbe1′ are almost equal due the closed loop. A bandgap voltage Vbg is determined according to the base-emitter voltage Vbe1 and Vbe2 (negative temperature coefficient) and thermal voltage VT (positive temperature coefficient) of the bipolar devices B11 and B12. With a proper ratio of the base-emitter voltage and the thermal voltage VT for each bipolar, the temperature dependence of the bandgap voltage Vbg is almost eliminated.\nHowever, because the area of the bipolar devices B11 and B12 is different, the current density in the bipolar devices B11 and B12 is also different, thus, the base-emitter voltage Vbe1 and Vbe2 of the bipolar devices B11 and B12 becomes unequal. Accordingly, the same current density in different areas of the bipolar devices B11 and B12 results in a current magnitude which is equal to ΔVbe/R13 and provides a positive temperature coefficient. The bandgap voltage Vbg is then determined according to the base-emitter voltage Vbe1 and a scale of ΔVbe and expressed as:\n  Vbg  =            Vbe      ⁢                          ⁢      1        +                            R          ⁢                                          ⁢          12                          R          ⁢                                          ⁢          13                    ×      Δ      ⁢                          ⁢      Vbe      \nTypically, the ratio of the resistors R12 and R11 is equal to about 10 to nullify the temperature coefficient.\nFIG. 2 depicts another conventional generating circuit of a bandgap voltage operating at low voltage, serving as a reference voltage for integrated circuits. Referring to FIG. 2, the generating circuit 2 mainly comprises an operational amplifier OPAMP, bipolar devices B21 and B22, and resistors R21 to R24. The values of the resistor R21 and R22 are the same. A bandgap voltage can be scaled freely by adjusting the ratio of the resistors and expressed as:\n  Vbg  =                    R        ⁢                                  ⁢        24                    R        ⁢                                  ⁢        22              =          ×              (                              Vbe            ⁢                                                  ⁢            1                    +                                                    R                ⁢                                                                  ⁢                22                                            R                ⁢                                                                  ⁢                23                                      ×            Δ            ⁢                                                  ⁢            Vbe                          )            \nAccording to the two generating circuits 1 and 2, the bandgap voltage of the generating circuit 2 is equal to R24/R22 of that of the generating circuit 1.\nIn above generating circuits 1 and 2, it is assumed that the operational amplifier OPAMP is ideal, meaning that the voltage of the two input terminals of the operational amplifier OPAMP is ideally equal. In practice, however, the voltage of two input terminals of the operational amplifier OPAMP is not equal, and there is a non-zero offset voltage Vos between the two input terminals. When the non-zero offset voltage Vos is considered, the generating circuit 1 of FIG. 1 is redrawn as in FIG. 3, and a bandgap voltage Vbg′ is expressed as:\n                              Vbg          ′                =                              Vbe            ⁢                                                  ⁢            1                    +                                                    R                ⁢                                                                  ⁢                12                                            R                ⁢                                                                  ⁢                13                                      ×            Δ            ⁢                                                  ⁢            Vbe                    +                                    (                              1                +                                                      R                    ⁢                                                                                  ⁢                    12                                                        R                    ⁢                                                                                  ⁢                    13                                                              )                        ×            Vos                                              (                  Equation          ⁢                                          ⁢          1                )            \nSimilarly, when the non-zero offset voltage Vos is considered, a bandgap voltage Vbg′ generated by the generating circuit 2 of FIG. 2 is expressed as:\n                              Vbg          ′                =                                  ⁢                                            R              ⁢                                                          ⁢              24                                      R              ⁢                                                          ⁢              22                                ×                      [                          (                                                Vbe                  ⁢                                                                          ⁢                  1                                +                                                                            R                      ⁢                                                                                          ⁢                      22                                                              R                      ⁢                                                                                          ⁢                      23                                                        ×                  Δ                  ⁢                                                                          ⁢                  Vbe                                +                                                      (                                          1                      +                                                                        R                          ⁢                                                                                                          ⁢                          22                                                                          R                          ⁢                                                                                                          ⁢                          23                                                                                      )                                    ×                  Vos                                            ]                                                          (                  Equation          ⁢                                          ⁢          2                )            \nAccording to Equation 1 and Equation 2, the non-zero offset voltage Vos of the operational amplifier OPAMP is amplified by (1+R22/R23). As described previously, in order to nullify temperature dependence of the bandgap voltage, the ratio of the resistors R2 and R1 is set to about 10. Thus, when the non-zero offset voltage Vos is considered, the bandgap voltage Vbg may drift from its ideal value by 10 times the non-zero offset voltage Vos. The amplified non-zero offset voltage Vos causes the reference voltage to vary greatly from chip to chip.\nIn order to reduce the variation of the reference voltage from chip to chip, several techniques are introduced. One method uses a trimming circuit. FIG. 4 depicts a trimming circuit. A bandgap voltage Vbg′ with a non-zero offset voltage Vos is amplified by an amplifier 40, and a trimmed output voltage Vref is selected from a resistor string 41. At the beginning of the trimming process, an ideal target voltage value is applied to an external tester and compared with an output voltage Vref of the trimming circuit. The output voltage Vref is selected by varying a final control code CODE40 which is determined according to the comparison result signal. When the output voltage Vref is equal to the ideal target voltage value, the final control code CODE40 is obtained. According to the final control code CODE40, the trimmed output voltage Vref is determined. The ideal target voltage value and the trimmed output voltage Vref can be fixed and remembered in an IC by fuses which can be programmed by laser or electronically. The trimming process, however, requires more time in the tester for calibration and fuses trimming.\nChopper stabilization is another method to reduce the variation in the reference voltage. U.S. Pat. No. 6,462,612 discloses a chopper stabilized bandgap reference circuit to cancel offset variation. Referring to FIG. 5, an input signal of an amplifier 51 is modulated by a high frequency modulator (MOD) 50. The modulated signal is amplified by the amplifier 51 and demodulated by a demodulator (DEMOD) 52. Since an offset voltage of the amplifier 51 is not modulated, the offset voltage is modulated to a high frequency in the demodulation process. The high frequency noise can be filtered by a low pass filter (LPF) 53. The chopper stabilized bandgap reference circuit is useful to reduce the variation caused by the offset voltage of the amplifier 51. The circuit requires a high frequency clock for modulation, however, which produces noise. Moreover, the low pass filter 53 occupies a large area."}
{"text": "1. Field of the Invention\nThe invention relates to a method for the discontinuous production of foam plastic parts, in particular of rigid polyurethane foam blocks, and to a specially designed mold for the execution of this method which is adapted to facilitate the production of homogeneous foamed plastic parts having substantially uniform properties throughout to include uniform compressive strength in the marginal regions.\n2. Description of the Prior Art\nVarious methods and associated molds are known in the art which attempt to achieve the desired characteristics of fixed shape with little scrap, homogeneity of structure, and uniformity of, among other things, marginal strength. The now commonly used methods for the discontinuous production of polyurethane foam in blocks provide either for closed foaming molds or molds with what are conventionally referred to as floating lids. Both methods, however, entail considerable disadvantages with respect to the obtainable foam quality.\nThe foam produced in closed molds has internal tensions and, due to the friction occurring at the mold walls during the foaming process, shows in the vertical marginal zones a drop in volumetric weight and in quality. Small deviations from the desired production values lead to uncompletely foamed upper edges or to tearing of the foam blocks, resulting in a high scrap rate. Due to the occurring foam pressures, closed block molds are relatively costly to produce.\nWhile the block foam made in molds with floating lid is produced without substantial pressure and shows no substantial tensions, the friction between the developing foam and the wall occurs fully in the same manner as in a closed mold. Here too, inhomogeneous qualities are caused in the vertical marginal zones due to the fact that the frictional forces increase with increasing foam level. The upper edges of the blocks are not square, thereby producing more scrap when the blocks are cut into panels."}
{"text": "1. Field of the Invention\nThis invention relates to an operational verification system for testing the operational capability of a missile approach warning system, such as is deployed on military aircraft, and methods for using that system. In particular this invention relates to an active coupler that creates and emits a test signal and which can be connected directly to electro-optical sensors deployed on aircraft to detect electro-optical signals emitted from missiles.\n2. Description of the Related Art\nToday's armed forces face increasing worldwide proliferation of missiles, including advanced infrared (IR) guided missiles, surface-to-air missiles, and air-to-air missiles. Political entities that once were confined to arms used in hand-to-hand combat have developed surface-to-air missiles. Missiles often attack without being visually observed, and can strike in a matter of seconds. Reliance upon and installation of missile warning systems is therefore increasing. Such systems are useful in multiple types of aircraft, and even in tanks to detect anti-tank missiles.\nAs a missile is essentially a variation on a rocket, each missile, as it travels towards its target, generates a plume, or exhaust trail, unique to that missile. Accurate plume analysis permits accurate identification of missiles and engagement in appropriate countermeasures. In analyzing the plume, the sensor counts photons in a sub-spectrum of interest of the ultraviolet or near-infrared spectrum. The sensor integrates the photon counts over small time intervals, which results in an optical power versus time signature, in watts over time, for each missile. This signature is compared to templates and the presence or absence of a match determined. The sensors reject man-made and natural clutter sources, and can detect missile plumes from a substantial distance.\nThe AAR-47 and 57 Missile Warning System/Missile Avoidance Warning System (“MAWS”), and other missile warning systems known in the art, are frequently used optical-based sensor systems deployed on fixed and rotary winged aircraft to detect short range missile launches. They consist of multiple electro-optical sensors, which read the signature and convey that data to an internal electronics control unit. The control unit processes the data, categorizes the missile, provides direction-of-arrival and elevation information, provides a warning, and directs countermeasures such as flares or jamming.\nSuch missile warning systems must be accurate; both early detection and a low rate of false alarms are required to protect pilots and their cargo without the stress of false alarms. The sensors must be very sensitive to the presence of plumes at a substantial distance, particularly when the aircraft is at a high altitude. Unfortunately, these missile warning systems rely on electronic and optical components that deteriorate with age and exposure to extreme environmental conditions such as those present at high altitudes and combat conditions (e.g., sand and salt water, in the case of aircraft launched from aircraft carriers). Sensitivity and accuracy of signature matching must be tested routinely in order to maintain optimal sensor performance. In addition, as new missiles with previously unknown signatures are developed, the sensors' accuracy in detecting those new signatures must be tested and evaluated.\nOne way of testing missile warning systems is by providing them with a signal in the same wavelength as a missile signature and evaluating the system's response. Such missile signature testing emits a signal in the appropriate wavelength, varying intensity (watts) and duration (time) of light presented to the sensor within the duration of the signal and between signals. This type of testing can be performed in the UV and infrared parts of the spectrum. Such a signal may be called a “waveform.” The sensor's response to the simulated signature is then analyzed for its ability to detect the waveform and direct appropriate countermeasures. More details about current waveform emission and sensor testing will be set forth below as its disadvantages are discussed in turn.\nCurrent testing systems present many problems. One concern is the level of extraneous light in many test situations, which weakens the test's accuracy. This is because the sensor counts both the photons in the environment, which are present in an uncontrolled and inconsistent amount, and the photons emitted by the tester. Welding, street lamps, and other sources can create environmental light in the spectrum generated by missile plumes and detected by missile warning systems. If light from these sources is detected by the sensor during a test, it will make the test less accurate because it is no longer based solely on the tester's calibrated emission. It is therefore desirable for a tester to be coupled to the sensor in a manner that excludes environmental light input and provides an isolated test environment. Herein, the term “coupler” refers to any mechanical enclosure that provides an isolated testing environment similar to that of a laboratory environment by preventing any undesired environmental effects.\nSome current testers have a coupler designed to block such ambient light. Such couplers are separate entities from the tester itself. They are generally tube-shaped and are held in place by the technician while he/she is simultaneously operating the sensor. This handheld manner of use introduces a great deal of human error, as the technician may easily inadvertently move the coupler, especially when tired or in harsh conditions as is probably the case in the context of an armed conflict or extended training exercise. Such movement may expose the sensor to ambient light while the tester is emitting its signal, thereby destroying the accuracy of that test and any grounds for comparison with other tests. In addition, the technician cannot walk away from the aircraft to attend to other tasks in testing the sensor, because the technician must hold the coupler.\nCurrent couplers also do not allow for accurate or repeatable positioning of the emitter and coupler. This may result in varying input angles that in turn adversely affect sensitivity testing data. That is, because the light shields of current testers are handheld and subject to inconsistent placement, the tester's signal is inconsistent in its intensity and directional approach relative to the sensor. Such inconsistency is unacceptable, as it is desirable to test the accuracy of a sensor's review of intensity and directional approach. It therefore remains desirable for a sensor to include a self-attaching coupler that completely and consistently blocks ambient light without relying on human positioning, and that standardizes the direction and intensity of the tester's signal.\nSubstantively, current tests are also far from comprehensive. Standard signature tests, called “Built-In-Tests” (“BITs”), are performed by a relatively small emitter that is part of the sensor itself as opposed to any exterior, more complex tester. BITs simply test whether the sensor does or does not detect any signal at all. Because the BITs do not generate a simulated signature or waveform, but rather a very simple “on/off” emission, they do not evaluate the sensor's accuracy in discerning between signatures. As such, BITs do not test a sensor's accuracy for different missile arsenals. This is particularly important as the field of missile technology advances and different countries generate different missiles. BITs also do not test the sensor's ability to accurately read the missile's angle of approach. Neither do they test the sensor's sensitivity, which is important as it reflects the sensor's ability to detect a plume at a substantial distance. BIT testing also does not test all quadrants of the sensor's field of view, leaving room for undetected inaccuracy. It is therefore desirable to expand standard sensor testing to test signature identification, all sensor quadrants, sensor sensitivity, and angle of approach.\nA more sophisticated current test is called Flight Line Test Set, or FLTS. FLTS uses low pressure mercury vapor lamps to generate a non-signature waveform: that is, the signal simply vacillates between low and high intensity, in the form of a light simply blinking on and off, rather than generating a signal that resembles an actual signature in its complexity. The sensor being tested thus receives not a simulation of a missile signature, but a relatively simpler “on-off” signal with no relation to any missile. This non-signature waveform does not adequately test the sensor's ability to read a complex signature. It also does not test the system's capability to detect and correctly recognize a specific missile threat, or to discern between different missiles. While the manufacturer of FLTS has limited FLTS to simple non-signature waveforms in the belief that actual signature testing is not necessary, such testing is believed important as the field of missile technology advances and different countries generate different missile arsenals.\nOne current improvement on FLTS, the Baringa, has a more accurate sensitivity test and can produce actual signatures. This system, too, is limited in its ability to simulate actual signatures, as its signals do not accurately represent missiles approaching from different directions. The generation of signatures representing approaches from different directions is performed by simply walking around the aircraft while transmitting a missile signature. This method is very inaccurate in representing a missile approaching from a specific direction. Sensors should be able to discern the direction from which a missile approaches, in order to provide the most useful information in the context of evasive or countermeasures (i.e., in which direction the plane should fly to avoid the missile, or in which direction antimissile projectiles should be launched).\nIt is therefore desirable for waveform testing to generate a simulated signature that can test the system's ability to discern between missiles and the direction from which those missiles approach. It is especially desirable for a tester to be equipped with software to produce signature simulations in accordance with set parameters, such as the actual signatures of missiles in an arsenal an aircraft may actually encounter in an upcoming mission. In this way, a fleet of planes or other vehicles headed for a conflict or warzone may be tested for their ability to detect the missiles prevalently used by enemy combatants in that area. In addition, it is desirable for these parameters to be reprogrammable, to adapt the tests to changing arsenals. In this way, planes or vehicles that move among or between conflicts or warzones over time may be tested for their ability to detect the most relevant arsenal.\nCurrent testers also do not provide any means for storing actual or simulated signature parameters, for reference in evaluating past tests or in creating and customizing future tests. Current testers cannot accomplish such storage of any information about the performed test. Testers with the capacity to internally generate waveforms and store test parameters and results are thus both self-sufficient and fully customizable. The customization according to stored parameters permits development and improvement of a threat library with signatures that simulate the plumes of the arsenal the aircraft may likely encounter. Another desirable use of memory is in the event of a sensor failure, to retroactively review what signatures were used in past tests of that sensor that permitted its failure to occur undetected.\nOne current tester, the Baringa, permits signature storage and generation of signatures according to stored parameters, but requires access to and exchange of information with a laboratory. Given the often remote locations in which aircraft sensors are tested, including aircraft carriers and deployments, this requirement of access to a laboratory hampers the ability to store signatures often when it is needed most, in battle or deployment. In addition, heightened security in those contexts often prohibits the exchange of data between a laboratory and the tester. It therefore remains desirable for a tester to be able to store signature data without requiring an external laboratory, in a simplified and portable theater.\nRegarding portability, it is particularly desirable for the generator and emitter to constitute one self-contained entity capable of enduring rugged conditions. It is also desirable that the tester not require calibration upon delivery to a test site. Portability is also enhanced by having a replaceable battery power source.\nCurrent testers only provide a very poor level of sensitivity testing. This is due in part to the inaccuracy inherent in the use of the handheld coupler, explained above, which lets in ambient light that the sensor may detect instead of or in addition to the faint sensitivity test signal. It is therefore desirable for a missile testing system to accurately test sensitivity. Sensors with good sensitivity are especially important as missiles become faster and it becomes more necessary to detect the missiles from a further distance in order to respond defensively.\nAnother disadvantage of FLTS is that each tester must be placed around three meters from the aircraft, with one operator per tester. Testing from a further distance requires another, separate set of long range testers; these multiple tester sets are expensive and tedious to transport and set up. It is therefore desirable for a missile tester to be able to test from any distance.\nFLTS emits its signal using low pressure mercury vapor lamps without a fluorescent coating. Such lamps present many problems. They use a great deal of power, which shortens the life of any associated battery power source. They are extremely heavy, which decreases the tester's portability and manipulability. In addition, such lamps require a high voltage, which creates a great deal of inefficiency in the form of heat when the lamps are turned on. Such lamps also waste operator time in that they require at least five minutes to “warm up” and stabilize before they can be used.\nMore specifically, mercury vapor lamps are not well suited to the specific task at hand of creating a variety of test signatures. It is not easy to change the intensity of a mercury vapor lamp's output, which stymies the instant purpose of creating different signatures based on varying intensity. These bulbs must be outfitted with bandpass filtering in order to produce definite wavelengths. This is problematic, as bandpass filtering does not in fact completely filter all undesired wavelengths, but simply attenuates those at the margins of desirability. Thus, testing with bulbs relying on bandpass filtering is not completely accurate. It is therefore desirable to use a sensor testing system that generates certain wavelengths without reliance upon bandpass filtering.\nAnother problem inherent in using mercury vapor lamps is that, in order to achieve different wavelengths, different specially designed bulbs must be built, purchased, and interchanged. This is costly, tedious, inefficient, and requires the risky manipulation of fragile tester components. It is therefore desirable for a sensor to be able to generate different wavelengths without bulb replacement.\nThe output between different sets of mercury vapor lamps is quite variable, as much as 30% between units; this detracts from the uniformity and ease of comparability between test results. It is desirable for a simulating tester to utilize bulbs other than low pressure mercury vapor lamps, in order to increase uniformity across testers.\nAn additional difficulty with current testers is that their filters are external and thereby prone to being damaged. Testers require filters in order to filter the signal from the strength at which it is generated to the strength appropriate for processing by the sensor. Some tests are unfiltered, while others are not; this depends in part on the tester's distance from the aircraft. Filters on current testers are external to the tester, such that they are easily scratched or shattered. The filter's external position is particularly risky due to the often harsh environments in which it is employed; aircraft are often in harsh environments wherein the filter may be damaged by sand, salt water, ice, or any other windblown particulate. In switching between filtered and unfiltered tests, current operators must manually remove the filter from the tester and stow it; this also introduces risk of damage to or loss of the filter. It is therefore desirable to have a filter internal to the tester which can be selectively engaged without separation from the tester."}
{"text": "1. Field of the Invention\nThe present invention relates generally to stimulating medical devices, and more particularly, to calibration of current sources and sinks in a stimulating medical device.\n2. Related Art\nThe delivery of electrical stimulation has become an established part of medical therapy. Numerous types of medical devices have components positioned on, or implantable in, a recipient's body in order to stimulate a recipient's tissue. Such devices are sometimes referred to herein as stimulating medical devices. Stimulating medical devices commonly include a plurality of electrodes that function as the interface between electronics of the device and the recipient's body tissue. In general terms, current is delivered to the recipient's tissue via the electrodes in order to evoke a response, such as a perception (e.g. for sound perception) or a function (e.g. for limb movement), in the recipient.\nFIG. 1 is a schematic diagram illustrating the delivery of current to tissue. In this illustration, an implantable stimulating medical device 110 comprises an implantable component 112, and a multi-electrode system in the form of two platinum electrodes 101. Each of the electrodes are connected to component 112 by insulated wires 104. Component 112 comprises a stimulating current source 102 that provides current to electrode 101A. The current passes through the recipient's tissue 100, including nerve cell 107, and returns to ground 114 within implantable component 112. The return of this current is shown schematically by arrow 103.\nAt the surface of platinum electrodes 101, chemical reactions take place, changing the electron current from the current source to an ion current 105 in the tissue. A charge 106 remains on the electrode surface, causing an increase in voltage in the tissue. Under normal conditions, these chemical reactions are reversible by a change in the direction of current. That is, a reversal in the direction of current will neutralize the increase in voltage. As such, it is common for the stimulation current to be delivered as biphasic pulses, in such a way that there is no net charge delivered to the tissue. A biphasic pulse includes a positive charge pulse followed by an equal negative charge pulse. In certain circumstances, the current level (amplitude) and periods of both the positive and negative pulses are substantially the same. In other circumstances, one of the pulses is applied over a longer or shorter period, but has lower or greater amplitude, respectively. However, in both circumstances, the total charge remaining in the tissue after delivery of both the positive and negative pulses is substantially zero.\nIn circumstances using biphasic pulses, if current is allowed to flow in one direction for too long, toxic products can escape and damage or destroy the surrounding tissue. Likewise, if the voltage between two electrodes is allowed to remain elevated for too long, toxic species are irreversibly generated. To ensure that stimulation remains safe, and that no toxic species escape, it is known that the DC and low-frequency (LF) states of the electrodes, sometimes referred to as the DC/LF voltages and the DC/LF currents, must remain within certain bounds. For a typical cochlear implant electrode having an area of about 0.25 mm2, these values are generally a few hundred milli-volts (mV), or tens of nano-amperes (nA). Additionally, the United States Federal Drug Administration (FDA) requires that the magnitude of the current through an electrode, during a 1 ms period, be below 100 nA. The use of charge neutralizing biphasic pulses helps ensure that these requirements are met, but charge errors occur in practice.\nIn certain stimulating medical devices, separate current source circuits and current sink circuits, referred to simply as current sources and current sinks, respectively, are configured to deliver or receive stimulating current. The sources and sinks each use a Digital-to-Analog Converter (DAC) to control the flow of current."}
{"text": "The subject matter herein relates generally to electrical connectors having electrical contacts that engage corresponding mating contacts of another electrical connector.\nElectrical connectors are used to transmit data and/or power in various industries. The electrical connectors are often configured to repeatedly engage and disengage complementary electrical connectors. The process of mating the electrical connectors may be referred to as a mating operation. For example, in a backplane communication system, a backplane circuit board has a header connector that is configured to mate with a receptacle connector. The receptacle connector is typically mounted to a daughter card. The header connector includes an array of electrical contacts (hereinafter referred to as “header contacts”), and the receptacle connector includes a complementary array of electrical contacts (hereinafter referred to as “receptacle contacts”). During the mating operation, the receptacle contacts engage and slide along the corresponding header contacts. The sliding engagement between the receptacle and header contacts may be referred to as wiping, because each receptacle contact wipes along an exterior surface of the corresponding header contact.\nElectrical contacts are often plated to enhance performance and/or durability of the electrical contact. Electrical contacts used to transmit data signals typically include one or more underlying materials and a plating disposed on the underlying materials. The outer plating is often a precious metal, such as gold, and is configured to impede corrosion while being sufficiently conductive for achieving the desired electrical performance. Plating material, however, that comprises gold or other precious metals can be expensive. But using less plating material to reduce costs can be problematic and may negatively affect electrical performance.\nAccordingly, a need remains for electrical contacts that require less plating material, such as gold and/or other precious metal, while achieving a sufficient level of electrical performance."}
{"text": "The present disclosure relates to a compound semiconductor layer formed of gallium nitride (GaN) or a nitride of gallium and other metal, and a method for forming the compound semiconductor layer. The present disclosure also relates to a method for preparing a substrate used in manufacturing of electronic or photo-electronic devices including the compound semiconductor layer. The present disclosure pertains to a technical field for forming a high quality compound semiconductor layer on a substrate, and more particularly, to a technical field for preparing a free standing compound semiconductor substrate by separating the substrate and the compound semiconductor layer.\nSemiconductor materials based on nitrides of Group III elements or Group V elements already hold important positions in electronic and photo-electronic fields, which will be important more and more. In fact, the nitride based semiconductor materials may be used in a wide range of fields from laser diodes (LD) to transistors operating at high frequency and high temperature. The nitride based semiconductor materials may also be used in ultraviolet photo-detectors, surface acoustic wave detectors and light emitting diodes.\nFor example, although gallium nitride is widely known for its usefulness in blue light emitting diodes and high frequency and high temperature transistors, it is also being extensively researched for use in microelectronic devices. As used herein, gallium nitride includes gallium nitride alloys such as aluminum gallium nitride (AlGaN), indium gallium nitride (InGaN) and aluminum indium gallium nitride (AlInGaN).\nTo grow a gallium nitride layer of low defect density is important in manufacturing gallium nitride microelectronic devices. A substrate on which gallium nitride is grown is known as one cause of the defect. However, it is difficult to prepare a gallium nitride substrate or a substrate for growing gallium nitride without defects. Typical methods such as Czochralski method where the crystal is grown from a melt cannot be used in producing a gallium nitride single crystal for the substrate because gallium nitride is difficult to melt. Surely, gallium nitride can be molten under ultrahigh pressure, however, this is currently unavailable for commercial use due to the low productivity.\nAccordingly, in such devices, the most frequently used for growing gallium nitride layer are heterogeneous substrates such as a sapphire substrate, a silicon carbide (SiC) substrate and a silicon substrate. However, because of lattice mismatch and thermal expansion coefficient difference between such substrate materials and gallium nitride, a large number of dislocations may be produced in the gallium nitride layer grown on the substrate, causing crack and bending of the gallium nitride layer. Therefore, a variety of buffer layers are often formed on the substrate before growing the gallium nitride layer thereon, or an epitaxial lateral overgrowth (ELO) method are used to reduce the dislocation generation.\nIn a typical ELO method, a stripe-shaped silicon dioxide (SiO2) mask is used to reduce stress caused by the lattice mismatch and the thermal expansion coefficient difference between the substrate and the gallium nitride layer. The typical ELO method will be described below with reference to FIG. 1, which is a cross-sectional view of the substrate on which the gallium nitride layer is grown according to the typical ELO method.\nIn the typical ELO method, the gallium nitride layer 2 is grown on the substrate 1 in a furnace, and then the substrate 1 is taken out of the furnace. The substrate is placed in a deposition apparatus so that silicon dioxide (SiO2) layer is deposited on the gallium nitride layer, and then the substrate 1 is taken out of the deposition apparatus. The silicon dioxide layer is patterned using a photolithography technique to form a silicon dioxide mask 3 on the gallium nitride layer, and then the substrate 1 is placed again in the furnace so that an ELO gallium nitride layer 4 is grown on the gallium nitride layer 2.\nA portion of the ELO gallium nitride layer 4 that is laterally grown over the silicon dioxide mask 3 has relatively high quality compared to the portion that is vertically grown. This is because defects such as dislocations cannot propagate through the laterally grown portion. Therefore, by forming a device in the portion of the ELO gallium nitride layer 4 that is laterally grown over the silicon dioxide mask 3, it is possible to obtain an excellent property.\nHowever, the ELO method requires the above described complex process such as an additional external process for forming the silicon dioxide mask, increasing process time and process cost. In addition, recently, as a plurality of silicon dioxide masks are used to improve and enlarge the function of the ELO, the number of the processes for forming the silicon dioxide mask and growing the gallium nitride layer is also increased correspondingly. Consequently, this may result in increased process cost, process complexity, time loss and economical loss, and thus result in decreased process yield.\nKorean Patent Laid-Open Publication No. 2004-0078208 discloses, instead of the ELO method, a method for preparing a gallium nitride substrate by forming grooves for reducing contact area between a sapphire substrate and a gallium nitride layer. According to the method, gallium nitride epitaxial layers are grown on an upper surface and a lower surface of the sapphire substrate, respectively. In specific, a first gallium nitride epitaxial layer is grown on the upper surface of the sapphire substrate in a furnace. The sapphire substrate is taken out of the furnace, turned upside down, and then placed in the furnace again so that a second gallium nitride epitaxial layer is grown on the other surface, i.e., the lower surface of the sapphire substrate. Next, a mask pattering is performed on the lower surface using a photolithography process, and the second gallium nitride epitaxial layer is etched to form a plurality of grooves. Thereafter, a laser beam is applied thereto to etch the portion of the first gallium nitride epitaxial layer corresponding to the portion of the plurality of grooves. As a result, the void-like grooves are formed on the upper surface of the sapphire. Using these grooves, a third gallium nitride epitaxial layer is grown.\nAccording to the above described method, the grooves prevent contact of the sapphire substrate and the gallium nitride layer for growing the third gallium nitride epitaxial layer. As such, it is possible to reduce dislocation generations in the gallium nitride layer over the grooves, and reduce crack and bending caused by the thermal expansion coefficient difference as temperature is decreased from a growth temperature to the room temperature. However, the method requires additional processes such as growing the gallium nitride layers on the upper surface and the lower surface of the sapphire substrate, respectively, performing the mask patterning on the gallium nitride layers, and applying the laser beam to form the grooves on the upper surface of the sapphire substrate. This may result in increased process time and process cost.\nIn addition, when the sapphire substrate is removed after growing the gallium nitride epitaxial layer on the sapphire substrate to use the gallium nitride epitaxial layer as a free standing gallium nitride substrate, the separation of the sapphire substrate from the gallium nitride epitaxial layer needs an additional process such as a laser lift-off. This may also increase the process cost, and decrease the process yield because heat applied to separate the sapphire substrate from the gallium nitride epitaxial layer may cause defects such as crack and bending in the gallium nitride epitaxial layer. If the substrate is formed of silicon, it may be easily removed by a polishing or a chemical etching. However, the silicon substrate also has a limitation that it is difficult to form a high quality gallium nitride epitaxial layer thereon.\nAs described above, although the mask patterning process for manufacturing the gallium nitride substrate of good quality is effective in reducing dislocations, it may increase process time and process cost."}
{"text": "1. Technical Field\nThe present invention relates in general to servo controls for a tape drive, and in particular to a technique for stabilizing a servo system for a cartridge tape drive that employs a tape on which servo and data tracks are arranged longitudinally.\n2. Description of the Related Art\nAround half a century has passed since the first open reel tape drives appeared, and today, tape drives continue to be widely used as mass backup storage for computers. Currently, small tape cartridge drives play a leading role and widely used for a variety of computers from low to high end. Generally, multiple data tracks are longitudinally arranged on a backup tape (normally, a magnetic tape) to record a large amount of data, and a tape head is precisely positioned on a target data track by servo control.\nThe timing based servo (TBS) system is a well known servo control system for a tape drive. According to this system, at least one specially patterned servo track is provided in parallel with data tracks, i.e., in the longitudinal direction of the tape, such that the timing of the servo position signal obtained by reading the servo pattern varies continuously as the tape head is moved across the width of the tape. To this end, the servo pattern is composed of magnetic transitions recorded in non-parallel directions, like a pattern composed of slash and backslash pairs, or a pattern composed of first and second chevron pairs with the first and second chevrons having different orientations. An exemplary TBS system is disclosed in U.S. Pat. Nos. 5,689,384 and 5,930,065 and U.S. patent application Ser. No. 09/370256 filed on Aug. 9, 1999.\nIn addition, several types of tape cartridge drives which implement the timing based servo, such as IBM 3570, are commercially available.\nAt present, a servo circuit for a tape cartridge drive (hereinafter referred to as xe2x80x9ctape drivexe2x80x9d or xe2x80x9cdrivexe2x80x9d) has filter coefficients which are determined based on a certain ideal model by taking mechanical and electrical margins into account. The mechanical margin relates, for example, to the tape feeding system and head system of the drive, tape distortion and servo pattern shifting, and the electrical margin relates, for example, to electrical disturbance, offset and computational error.\nHowever, even if the filter coefficients are determined by taking these margins into account, there are some drive and cartridge combinations which lack servo stability. For example, some drive and cartridge combinations may have an impulse-like spectrum with an abnormal magnitude at a certain frequency. Current servo systems may have a setup which amplifies this frequency component so that the servo can not follow the tape, which eventually results in a failure of tape writing. Although it is desirable to obtain an optimum set of filter coefficients for all drive and cartridge combinations, it is practically impossible to obtain such optimum filter coefficients since the frequency component which causes instability is not fixed and varies depending on the drive and cartridge combination.\nIt is, therefore, an object of the present invention to stabilize a servo system by eliminating the cause of instability due to a particular tape drive and cartridge combination.\nIt is another object of the present invention to provide a method and system for stabilizing a servo system by re-setting filter coefficients each time a tape cartridge is inserted into a tape drive.\nAccording to a first aspect of the present invention; in a tape drive system performing read/write operations for a tape by inserting a tape cartridge into a tape drive, where the tape cartridge has a tape on which servo and data tracks are arranged longitudinally and a predetermined servo pattern is recorded on each of said servo tracks, a method for stabilizing a servo system for the tape drive is provided which comprises the steps of re-setting coefficients of a low-pass filter for servo control of the tape drive, each time the tape cartridge is inserted into the tape drive, in response to the inserted tape cartridge.\nMore specifically, the method according to the first aspect comprises the steps of reading a servo pattern recorded on the tape in said tape cartridge inserted into the tape drive, performing a frequency analysis of the servo pattern, detecting, from frequency components obtained by the frequency analysis, a singularity that exceeds a predetermined spectrum range, and setting said coefficients of the low-pass filter so as to cancel a power spectrum of the singularity.\nAccording to a second aspect of the invention, a tape drive is provided into which a tape cartridge is inserted, the tape cartridge having a tape on which servo and data tracks are arranged longitudinally and a predetermined servo pattern is recorded on each of the servo tracks, and which comprises a servo control circuit including a low-pass filter for servo signals, and a circuit for re-setting coefficients of the low-pass filter each time a tape cartridge is inserted, in response to the inserted tape cartridge.\nPreferably, the circuit for re-setting the coefficients of the low-pass filter comprises a circuit for reading a servo pattern recorded on the tape in the tape cartridge inserted into the tape drive, a circuit for performing a frequency analysis of the servo pattern, a circuit for detecting, from frequency components obtained by the frequency analysis, a singularity that exceeds a predetermined spectrum range, and a circuit for setting the coefficients of the low-pass filter so as to cancel a power spectrum of the singularity."}
{"text": "This specification relates to processing computational graphs representing neural networks.\nNeural networks are machine learning models that employ one or more layers of models to generate an output, e.g., one or more classifications, for a received input. Some neural networks include one or more hidden layers in addition to an output layer. The output of each hidden layer is used as input to the next layer in the network, i.e., the next hidden layer or the output layer of the network. Each layer of the network generates an output from a received input in accordance with current values of a respective set of parameters for the layer.\nThe layers of a neural network can be processed by an individual device. The device can have a processor that performs operations, e.g., generating outputs at a layer from inputs, and stores outputs from the operations in memory. Due to the large number and size of operations generally required to generate the outputs in the neural network, one device can take a significant amount of time to process the layers of the neural network."}
{"text": "Technical Field\nCertain embodiments of the present invention relate to a forming system and a forming method adapted to form a metal pipe.\nDescription of Related Art\nForming systems that perform forming by expansion with the supply of a fluid into a metal pipe material have been known. For example, a forming system disclosed in the related art is provided with a pair of upper and lower dies, a holding unit that holds a metal pipe material between the upper die and the lower die, and a fluid supply unit that supplies a fluid into the metal pipe material held in the holding unit. In this forming device, the metal pipe material is expanded by supplying a fluid into the metal pipe material held between the upper die and the lower die, and thus can be formed into a shape corresponding to a shape of the die. Such a forming method is referred to as hydroforming."}
{"text": "This invention relates to a simple method for enhancing the weak colors displayed by certain iridescent pigments, so as to provide more vivid colors which are useful in printing, copying, and recording systems. This invention further relates to a novel method for recording indicia by means of iridescent colors. This invention still further relates to a novel method for forming intensely colored indicia from weakly colored recording fluids. This invention still further relates to a new series of novel printing inks.\nIn recent years there has appeared a new type of nacreous pigment which exhibits a weak iridescence when dispersed in liquid. These pigments are known commercially as \"iridescent pigments.\" From an optics view point, they comprise a base substrate of essentially transparent (colorless) plate-like particles, with at least one overcoating of an essentially transparent material having a refractive index greater than that of the base substrate. They are prepared commercially by overcoating a micaceous substrate with an extremely thin film of a material such as titanium dioxide or zirconium dioxide, having a refractive index which is substantially higher than that of the base substrate and which is usually greater than about 2.0. Such pigments are disclosed in U.S. Pat. Nos. 3,071,482 and 3,087,828 and are available from Mearl Corp., 41 East 42nd St., New York, N.Y. 10017, under the name \"Flamenco Nacreous Pigments.\"\nIn the course of another investigation, the present applicant examined these iridescent pigments for possible use in a recording paper. They displayed very weak colors which gave almost no contrast with normal backgrounds and, therefore, they had no practical value for that purpose. As graphite is strongly colored and has a plate-like structure similar to mica, an attempt was made to deposit an iridescent interference layer, the so-called Nobili rings, on graphite flakes, but the Nobili process did not work on graphite. (The Nobili rings and process are discussed in David Fishlock, Metal Colouring, Robert Draper, Ltd, Teddington, 1962.)"}
{"text": "Field of the Invention\nThe present invention relates to the technical field of display technology, and specifically to a mounting structure and a wall-mounted display device.\nDescription of the Related Art\nSince manufacturing tolerances exist in display screens and mounting brackets, and assembly tolerances exist between display screens and mounting brackets, relative large gaps or overlaps exist between assembled displays if no adjusting means is provided. One recently used adjustable bracket is shown in FIG. 1, which comprises a fixed bracket supporting a display device 1, the fixed bracket is provided with two wall-mounting members 3 mounted onto a mounting beam 5. Each wall-mounting member 3 is provided with a vertical adjusting bolt 2 and a horizontal adjusting bolt 4 for adjusting the member vertically and horizontally, respectively. The vertical adjusting bolt 2 and the horizontal adjusting bolt 4 are connected to the wall-mounting member 3 through threads, and only contact the mounting beam 5. Therefore, when threaded toward inside of the wall-mounting member 3, the adjusting bolts press against the mounting beam 5, so as to change gaps between the wall-mounting member 3 and the mounting beam 5 in the vertical direction and in the horizontal direction, and thereby an adjustment of position of the display device 1 is achieved. However, when threaded outwards from the wall-mounting member 3, the horizontal adjusting bolt 4 is disengaged from the mounting beam 5, such that it cannot move the wall-mounting member 3 rightward in relation to the mounting beam 5, as shown in FIG. 1. Therefore, the adjustable bracket of this type has a limitation.\nTherefore, a wall-mounting structure is needed to adjust a position of a display device in horizontal and vertical directions."}
{"text": "1. Field of the Invention\nThe present invention relates to a security policy management system and method for deriving a security policy from setting details of a security device that is a component of an information system.\n2. Description of the Related Art\nWith the development of the information technology, information security is acquiring growing importance. Today's organizations including enterprises and other institutions are required to have a security policy established from security setting details in their network systems, so that the system administrators and security managers can understand the security policy.\nOne technique for grasping a security policy is disclosed as a status grasping method in Japanese laid-open patent publication No. 2003-203140 (JP, P2003-203140A). The disclosed status grasping method comprises the steps of asking questions of the members of the target organization, grasping a security situation of the information system of the target organization based on the answers to the questions, grasping a security situation of the information system of the target organization based on the results of an investigation performed by an investigation tool, and establishing a security policy by integrating the information obtained from the above two grasping steps. Publication JP, P2003-203140A also discloses an integrating apparatus for use in carrying out the above status grasping method. The integrating apparatus compares the status of the information system according to the answers to the questions and the status of the information system according to the investigation tool. Portions of these two statuses that match each other are referred to as matched portions, and portions of these two statuses that do not match each other are referred to as unmatched portions. The integrating apparatus separates the unmatched portions from the matched portions, takes out the unmatched portions, and displays the unmatched portions. The integrating apparatus then prompts the user as the system administrator to enter the user's selection with respect to the unmatched portions, displays the result selected by the user, and combines the matched portions and the result selected by the user.\nSince the method and apparatus disclosed in JP, P2003-203140A uses not only questions for the members of the organization but also the investigation tool, it is not necessary to ask questions about items that can be investigated with the investigation tool. Furthermore, as the unmatched portions are displayed to prompt the user to make a selection, the user can make a suitable selection from the displayed details, allowing a more appropriate security policy to be established.\nThe conventional method and apparatus depends upon the experiences and knowledge of the security manager, and constructs a security policy based on the answers to the questions that the members of the organization were asked. Consequently, the conventional method and apparatus suffer problems, to be described below, that are caused by the fact that a security policy is constructed by manual intervention.\nFirst, a large expenditure of labor and time is required to construct and grasp a security policy. A security policy needs to be constructed for each individual organization, and the security manager has to create a security policy from scratch based on his or her experiences and knowledge, and to gather the answers to the questions that the members of the organization were asked. Setting information of existing security devices in the information system of the organization is described according to a format inherent in each security device. There is a tacit understanding with respect to each security device that the description of some information is omitted. A security policy derived from such setting information lacks integrity and is difficult to grasp unless the security policy is manually corrected and adjusted. As a result, a large expenditure of labor and time is required to construct and grasp a security policy.\nA second problem is that there is a high possibility of errors and omissions occurring in constructing a security policy and grasping a security status. The reasons for such a high possibility of errors and omissions are that the manual intervention in the process of constructing a security policy is likely to cause errors and omissions due to assumptions and illusions that the security manager and the respondents to the questions may have.\nAccording to JP, P2003-203140A, a scanner is disclosed as the investigation tool for investigating the security status of the information system. However, such an investigation tool may possibly collect erroneous information which is different from the setting details that are actually set in a security device to be investigated.\nSecurity policy management systems for deriving a security policy from setting details of a security device that is a component of an information system are to meet the following demands:\nWhen a security policy is constructed, a security policy management should preferably be able to present the details of the security policy in an understandable way to the manager. A security policy management system should preferably be able to compare the security policies of a plurality of systems, and to compare not only the security policies of a plurality of systems, but also security policies with each other, e.g., the security policy of one system and a predetermined reference security policy with each other.\nA security policy management system should preferably be able to analyze a plurality of security policies of the same type in an integrated manner. For example, if there are a plurality of packet filtering devices, then a security policy management system should preferably be able to analyze security policies grasped respectively from the packet filtering devices and specify packets that can pass through all the packet filtering devices.\nA security policy management system should preferably be able to interlink security policies or the like having settings for different functions and grasp the relation between those security policies. For example, a security policy management system should preferably be able to interlink a security policy having settings for a packet filtering function and the results of an analysis of settings for an intrusion detecting function, and to determine whether there is an unmatch between the settings or not."}
{"text": "The present invention relates generally to engineering problem solving and design information processing systems and more particularly to computer based systems for aiding engineers, scientists and the like to have a greater understanding of the products, processes, or machines they wish to improve and the technical problems related thereto that they wish to solve.\nGreat advancements have been made in the fields commonly known as computer aided design (CAD) and computer aided engineering (CAE). These computer based systems enable the designer to create detailed images and print-outs of the product, process or machine he/she is designing or improving. With CAD, the designer can try many new designs or modifications of subsystems and components quickly and view the modified products immediately on the monitor or print-out. The CAD system also generates virtual 3-D images of the product or machine, enables in-space rotation of the product image and zoom through the product image interior.\nAlthough CAD systems are a great designer tool for trying design changes quickly, they do not otherwise aid the designer in the evaluation and solving of technical engineering problems or conceiving new products or processes that provide new functional performance or the same functional performance with completely different engineering approaches. Accordingly, there has arisen a relatively new area of computer based engineering tools known as concept engineering computer based systems. These systems serve to increase the designer's inventive and creative abilities in solving engineering and scientific operational or functional problems and, in the course of such problem solving, induce the designer to invent new structural and functional concepts applicable to his/her design goals.\nOne such concept engineering computer based system is the TECHOPTIMIZER.TM. software sold by Invention Machine Corporation of Boston, Mass. and described in pending U.S. patent application Ser. No. 08/822,314, filed Nov. 12, 1996 that comprises a knowledge and logic based information processing system for generating conceptual engineering system definition and problem analyses. This system automates the process of originating the statement of the most important technical problems for elimination of components or harmful inter-actions between components of the object system being analyzed or redesigned.\nAnother concept engineering computer based system is the Invention Machine.TM. LAB.TM. Software sold by Invention Machine Corporation of Boston, Mass., that comprises a knowledge and logic based system that generates concepts and recommendations for solving engineering problems at the conceptual level. Various inventive rules or procedures are included and certain ones are selected and presented to the user for consideration in solving the user's current session problem. This system applies to all fields of physical science and aids the engineer by solving engineering contradictions to reduce the tendency of user applied engineering trade-offs. The system includes a large data base of physical, geometric, and chemical effects used in the past to solve other engineering problems. Selected ones of the effects are presented to the user for consideration by the user as potential solutions to his/her current problem session. The system also includes a technology evolution and prediction capability that aids the user in understanding the dynamics of his/her product evolution and the logical next or future generation of the product or its function. This stimulates the user to think forward and extrapolate the dynamics of the technology life cycle and originate the next generation of technology.\nAlthough the above mentioned computer based concept engineering systems have experienced much acceptance by the technical community, there is still a need for alternate methods of concept engineering that enable the user to rapidly and comprehensively understand the nature of the technical result he/she wishes to achieve and the structure and function of various engineering devices that, if linked together, can produce the desired technical result.\nThere is also a need in concept engineering and design systems for the computer based system to rapidly, effectively and visually present to the user the capability of searching for the precise technical result he/she desires and displaying to the user in brief but comprehensive form an engineering system or device, what the physical effect is, and how the device can produce the desired technical result. The problem here is generating for the user information in such a way that enables the user to obtain and understand substantially immediately, e.g. within one or two minutes, the technical problem the user must address to achieve the desired technical result.\nIn addition, there is a need to present to the user/designer for consideration various subsystems that either control or affect the parametric environment for the engineering device under consideration or that can function as a series of input-output stages for the system representation that produces the desired technical result."}
{"text": "1. Field of the Disclosure\nThis application generally relates to urinary catheters, assemblies including urinary catheters, and medical uses thereof.\n2. Background of the Disclosure\nSome patients require urinary catheters to effectively and hygienically void their bladder, such as those patients who suffer from urinary incontinence or other medical disorders.\nUrinary catheters are inserted into the bladder through the urethra in order to open the urinary sphincter and drain urine into an appropriate receptacle. A typical urinary catheter patient might void their bladder every two to four hours. In developed countries, urinary catheters are disposable (that is, they are not reused), with the effect of requiring the use of multiple catheters per day.\nIt sometimes occurs that urinary catheters for patient use are insufficiently compact, or otherwise involve relatively bulky equipment. For example, in current use, many catheters are packaged in an elongated form, with the effect that they are cumbersome for both storage and transport. This can have the effect that the urinary catheter, or its presence on the patient, is readily apparent, or is not easily conveyed by the patient or in the patient's clothing when in use, or is not easily concealed or otherwise attached to the patient's clothing when in use.\nSimilarly, it sometimes occurs that using a urinary catheter has the substantial possibility of social embarrassment for the patient due to the possibility of the catheter, or any associated medical equipment, or the fact of the patient's need for the urinary catheter, becoming observed in a social setting. This can also occur due to the urinary catheter, or its presence on the patient, becoming apparent, or not being well concealed by the patient's clothing. More recently, smaller and more discreet catheters have been developed, such as the “SpeediCath™ Compact” catheter by Coloplast, Inc., currently available for female users in the United States, and available for male and female users in Europe.\nIt sometimes occurs that it is difficult to handle lubricated catheters, as they are often designed to be slippery, so as to facilitate easy insertion into the urethra. For example, in current use, many catheters are either pre-lubricated with a water activated coating, or are designed to be lubricated with a gel after opening, in either case in order to make it easier for the urinary catheter to be inserted into the urethra. Moreover, these lubricating mechanisms can be inconvenient, as users must either deal with packaging fluid, or carry around a lubricant along with the urinary catheter.\nIt also sometimes occurs that using a urinary catheter can be unhygienic due to the possibility of touching the catheter with an unsanitary object, such as either the patient's hand, or a portion of the patient's body other than for insertion into the urethra. This can have more than one untoward effect. A first effect can be that use of an unhygienic urinary catheter might lead to a urinary tract infection, or other untoward medical condition. For example, when handling the catheter directly, users must generally ensure that their hands and the opening of the urethra are clean, in order to avoid contaminating the catheter and risking a urinary tract infection. A second effect can be that the urinary catheter must be disposed of after use, and replaced with another such catheter, with the concomitant difficulty of needing additional urinary catheters, as well as some form of storage for used catheters.\nIt also sometimes occurs that patients using a urinary catheter are reluctant to detach the catheter from its packaging for voiding their urine. For example, in current use, urinary catheters can be conspicuous and difficult to use, with the effect that self-conscious patients might be inclined to refrain from using them sufficiently often, particularly in certain social settings, with the concomitant possibility of bladder problems, kidney problems, or other medical conditions."}
{"text": "The present invention generally relates to projection cockpit displays, and even more particularly relates to such displays having multiple tiled display surfaces.\nIn the past, designers of avionics displays have endeavored to provide projection displays with ever-increasing image quality, reliability, and economic efficiency. Projection displays based on liquid crystal display (LCD) technology generally use a single, linear polarization state for light incident upon the display element. Typically, this is accomplished by the use of a linearly polarizing filter, prism, or reflector disposed in front of the LCD device. Linearly polarizing filters, prisms, and reflectors are all well known in the art. These polarizing devices typically transmit one linear polarization state, while a second orthogonal polarization state is either reflected into a different path, or absorbed within the polarizing element. In either case, the state which is not transmitted becomes unusable by the display. This results in a significant (typically greater than 50 percent) attenuation of the light flux which is usable by the display.\nIn the past, display designers have proposed methods and devices to capture and reuse the otherwise discarded orthogonally polarized light and direct it into the display device, along with the polarized light originally produced from the incident light. These devices have been called xe2x80x9cpolarization recovery devicesxe2x80x9d, which are typically an array of prisms and half-wave retarders which divide the light into two polarization states, rotate one polarization state by 90 degrees, then recombine the now parallel polarization states. These xe2x80x9cpolarization recovery devicesxe2x80x9d have found limited acceptance in the LCD projector industry because they are quite expensive to fabricate, and they can also result in significant other performance deficiencies.\nConsequently, there exists a need for improvement in illumination of tiled projection displays.\nIt is an object of the present invention to provide tiled projection displays having a superior image luminance characteristic.\nIt is a feature of the present invention to utilize multiple tiled viewing surfaces which receive illumination through the same front-end polarizer.\nIt is another feature of the present invention to utilize a polarizing prism or reflector as the front-end polarizer.\nIt is an advantage of the present invention to optically divide light from a single light source and use the divided light through multiple LCDs.\nIt is another feature of the present invention to include a half wave retarder to rotate the polarization state of light directed to an adjacent LCD.\nIt is another advantage of the present invention to increase commonality of component parts by having the same polarization state for all light incident on the LCDs in both optical channels in pairs of tiled projection displays.\nThe present invention is an apparatus and method for illuminating tiled projection displays designed to satisfy the aforementioned needs, provide the previously stated objects, include the above-listed features, and achieve the already articulated advantages. The present invention is carried out in an efficient manner that wastes less light flux than is typically done when an absorbing polarization filter is used to provide linearly polarized light.\nAccordingly, the present invention is a tiled projection display system which utilizes a tiled multi-viewing surface display where multiple viewing segments are illuminated by light, which has been conditioned by the same polarizing device prior to its passing through separate liquid crystal display (LCD) devices."}
{"text": "1. Field of the Invention\nThe present invention relates generally to systems and methods for motivating employees, and more particularly to advantageous aspects of systems and methods for providing performance feedback to a cashier at a point-of-sale (POS) terminal.\n2. Description of the Prior Art\nIt is known that one way to motivate employees to perform their tasks as efficiently as possible is to provide them with ongoing feedback as to their level of performance. In a retail store, one important task having a direct impact on customer satisfaction is the speed and accuracy with which the checkout process is accomplished. However, cashiers using POS terminals currently do not get immediate feedback on their performance. Thus, providing cashiers with this feedback may motivate them to work faster."}
{"text": "1) Technical field of the Invention\nThe present invention relates to a power semiconductor device, and in particular, relates to the power semiconductor device with improved reliability for a thermal shock test.\n2) Description of Related Arts\nA power semiconductor device such as a power module includes, in general, an insulating substrate and a heat radiator (heat sink) bonded thereon via an adhesive layer of electrically conductive material such as a solder layer. Mounted on the substrate are a plurality of semiconductor chips such as an insulating gate bipolar transistor (IGBT) and a free wheel diode (FWD).\nThe power semiconductor device is often exposed to the thermal shock due to a varying ambient temperature and heat generated from the semiconductor chips in operation. Difference of linear expansion coefficients (linear expansivity) between the insulating substrate and the heat radiator causes a substantial amount of strain (stress) within the intervening solder layer, in which the solder cracking is occurred. When heated, the heat radiator made of metal such as copper expands relative to the insulating substrate from the middle portion towards the circumference thereof, and when cooled, the heat radiator shrinks relative to the insulating substrate from the circumference towards the middle portion thereof. Therefore, the stress of the solder layer adjacent to the circumference is much greater than that adjacent to the middle portion. Thus, the solder cracking is firstly appeared on the circumference and extends towards the middle portion of the solder layer as the power semiconductor device is repeatedly exposed to a number of the thermal shocks.\nThe solder cracking of the solder layer may arise various problems, in particular, when the solder cracking extends through the area beneath the semiconductor chips generating a substantial amount of heat, it likely blocks the heat conducting from the semiconductor chips to the heat radiator. Thus, the solder cracking may cause the semiconductor chips to be overheated, and hence to malfunction. In order to improve the reliability for the thermal shocks, it is required to prevent and/or delay the extension of the solder cracking from the circumference towards the middle portion of the solder layer so that no solder cracking is generated especially beneath the semiconductor chips.\nMany approaches to prevent and/or delay the extension of the solder cracking has been incorporated in the conventional power semiconductor device. For example, the first approach is to design the intervening solder layer with increased thickness to absorb the stress therein. The second one is to design the insulating substrate having a round corner with the increased radius of curvature so that the distance from the corner to the middle portion is reduced.\nThe third one is to arrange a plurality of different semiconductor chips on the insulating substrate in a symmetric manner so that the deviation in the thermal distribution of the insulating substrate heated by each semiconductor chip is minimized. The fourth one is to develop a new material of the insulating substrate for reducing the stress in the intervening solder layer.\nAlso, according to the fifth approach, the semiconductor chips are provided on the insulating substrate as far away as possible from the circumference, thereby to prevent the extension of the solder cracking through the area beneath the semiconductor chips.\nOn the other hand, according to the semiconductor device disclosed in JPA 10-50928, a plurality of bosses are provided between the insulating substrate and the heat radiator in order to control the intervening solder layer to be thicker than a predetermined thickness.\nAlso, according to the heat sink of the semiconductor device disclosed in JPA 10-189845, a peripheral trench is provided on the heat sink at the position corresponding to the circumference of the insulating substrate, which is filled up with the solder. This effectively increases the thickness of the solder layer so that the strain in the circumference of the solder layer is absorbed.\nHowever, there are several drawbacks with the above-mentioned approaches. That is, the intervening solder layer with the increased thickness causes a poor thermal conductivity of heat from the insulating substrate to the heat radiator via the intervening solder layer. Also, the package or size of the semiconductor device is increased because of the insulating substrate having the round corner (the second approach). Also, much greater space or area for mounting the semiconductor chips is required if the semiconductor chips are arranged on the insulating substrate in a symmetric manner (the third approach), or at positions as far away as possible from the circumference thereof (the fifth approach). Further, even though the new material of the insulating substrate has been developed (the fourth approach), it still needs time and efforts to adapt the material to the actual products.\nFurthermore, the bosses indicated in JPA 10-50928 may be used to control the intervening solder layer to be thicker than a predetermined thickness. However, it can hardly controlled such that the thickness of the solder layer is uniform across the solder layer in a precise manner. When the insulating substrate is inclined relative to the heat sink in a certain direction, any one of the corners of the intervening solder layer has the thinner portion, so that the solder cracking extends from the thinner portion to the area beneath the semiconductor chips.\nIn addition, the trench disclosed in JPA 10-189845 may increase the thickness of the solder layer at the circumference, yet, the semiconductor device has to be mounted on the heat sink in a flat or parallel manner.\nTo address the aforementioned problems, one of the aspects of the present invention is to provide a power semiconductor device including a heat radiator having a principal surface and an insulating substrate bonded on the principal surface of the heat radiator via a first solder layer. The power semiconductor device also includes at least one semiconductor chip mounted on the insulating substrate via a second solder layer. The insulating substrate has a thin-layer and thick-layer edges, and is bonded on the principal surface of the heat radiator so that the first solder layer has a thickness thinner towards a direction from the thin-layer edge to the thick-layer edge (T1 greater than T2) . Also, the semiconductor chip is mounted on the insulating substrate so that a first distance between the thick-layer edge and the semiconductor chip is less than a second distance between the thin-layer edge and the semiconductor chip (L1 less than L2).\nFurther scope of applicability of the present invention will become apparent from the detailed description given herein. However it should be understood that the detailed description and specific examples, while indicating preferred embodiments of the invention, are given by way of illustration only, since various changes and modifications within the sprit and scope of the invention will become apparent to those skilled in the art from this detailed description."}
{"text": "The first function of the epidermis is to form a barrier between the external environment and the internal medium. It is the outermost layer of the epidermis, the stratum corneum, which carries out this task. It is formed of keratinocytes at the final stage of their differentiation: corneocytes, which are cemented to one another by a thick intercellular cement which is both flexible and impermeable. This physical barrier which is the skin makes it possible for the human body to protect itself against numerous types of attack. These attacks can have different intrinsic origins, such as chronological aging or else biochemical changes which take place during states of fatigue, stress or hormonal changes such as those during pregnancy, etc. Other attacks are extrinsic in origin, such as pollution, the sun, disease, etc. In response to these attacks, the appearance of the skin is changed and the appearance of wrinkles and fine lines, areas of hyper- or hypopigmentation, dryness or even dehydration of the skin, a thinning of the epidermis, elastosis, imperfections, aged areas, etc. is observed. These changes are caused by the alteration to the functions of cellular renewal, cellular cohesion and synthesis of collagen, elastin and other proteins, and ultimately lead to a decrease in the protective barrier qualities of the skin and to a less attractive appearance thereof.\nLike all the other organs, the skin is subjected to the influence of periodic variations. It represents the head office of a coherent organization of circadian and nycthemeral rhythms which modulate various, usually large-scale, biological cycles (F. Henry et al., Rev Med Liege; 57:10:661-665). In humans, these cutaneous rhythms suggest that, during the day, the skin promoted various protective functions with regard to the environment. During the evening and night, it promotes its cellular renewal and various metabolic synthesis processes.\nIn a subject in good health, where the organism lives in harmony with its environment, the biological rhythms are synchronised. By contrast, disturbances to the biological rhythms may appear in a certain number of conditions called ‘desynchronization’ (Reinberg and Touitou, 1996). Desynchronization is a state in which two (or more) rhythmic variables, which were previously synchronised, no longer have the same frequency and/or acrophase relations and exhibit temporal relations which are different to usual relations. Desynchronization may be of external origin, and thus depends on environmental changes and occurs, for example, during a transmeridian flight passing through five time zones (‘jet lag’) or when working through the night. Desynchronization on internal origin does not depend of environmental factors. It occurs with aging or with a certain number of illnesses, such as depression or some cancers. These problems of desynchronization can be corrected by various treatments: for example the administration of strong light in order to treat seasonal depression, or else the administration of melatonin (Dijk et coll., 1995; Eastman et Miescke, 1990; Palm et coll., 1992; Schochat et coll., 1998; Touitou et coll., 1998). In cosmetics, it is well known that the efficacy of a treatment can be optimized in accordance with the time of day at which it is administered. For example, day creams make it possible to protect the skin against external attacks during the day, whereas night creams enable the skin to repair any damage incurred during the day by increasing cellular renewal and metabolism. Furthermore, it has been demonstrated that aging is accompanied by a change in the biological rhythms, with a decrease in amplitude and a tendency towards phase lead (Weinert D., Chronobiol. Int. 2000; 17:261-83). Moreover, it has been shown that the circadian clock was altered by this same aging process. However, no treatment currently makes it possible for the skin to ‘re-start’ or else revive its circadian cycle, or to resynchronize its biological clock. Such a treatment would thus make it possible to help the skin to recover from a time difference, night-time work, or else to fight against signs caused by aging.\nSurprisingly, the applicant has discovered that peptide compounds of the following general formula R1-(AA)n-X1-Ser-Thr-Pro-X2-(AA)p-R2 have the property of restoring the circadian rhythm and resynchronising the biological clock of skin cells. No document of the prior art describes such peptide compounds in order to obtain such effects. Furthermore, these peptide compounds are characterized by the fact that they are activator agents, either directly or indirectly, of the Clock (circadian locomotor output cycles kaput) protein involved in the regulation of the circadian cycle. Consequently, the present invention relates to peptide compounds which are Clock activators and to their use in cosmetic and pharmaceutical compositions in order to restore the circadian rhythm and resynchronize the biological clock of skin cells or else prevent or correct the signs caused by aging.\nIt has long been known that the circadian clock is controlled by a negative regulation loop involving a set of genes, in particular the Per-1 (period), Clock and BMAL-1 (brain and muscle ARNt-like protein) genes. Previously, in the prior art, numerous applications have already been proposed with regard to the use of nucleotides and/or proteins produced by Clock, Per-1 or BMAL-1 genes. For example, see U.S. Pat. No. 6,291,429 which describes the use of the product of the Clock gene to resynchronize the sleep cycle or physiological or endocrinal processes, or to resynchronize the body after jet lag, etc. However, no document of the prior art describes the use of specific Clock activator peptides in order to restore the skin's internal biological clock."}
{"text": "Currently, the number of various crimes such as robbery of articles placed inside cars and illegal entry of another's residence has been increasing, and the consciousness of people for crime prevention has been growing. To meet the demand generated with the increase in consciousness for crime prevention, a variety of crime-prevention goods have been developed.\nFor example, a security mechanism control system which locks and unlocks doors of a vehicle has been developed as an example of a security mechanism control system for a vehicle using radio waves. For example, according to security mechanism control systems disclosed in Patent Reference No. 1 (JP-A-2003-27790, publication date: Jan. 29, 2003), Patent Reference No. 2 (JP-A-2003-27791, publication date: Jan. 29, 2003), and Patent Reference No. 3 (JP-A-2003-301638, publication date: Oct. 24, 2003) or other references, a user can automatically lock or unlock a vehicle only by carrying a portable device and easily check that the vehicle has been locked. In these security mechanism control systems, the portable device transmits signals and a main device receives the signals. Thus, the user carrying the portable device can automatically unlock the vehicle only by approaching the vehicle and lock the vehicle by moving away from the vehicle.\nIn such a conventional type of security mechanism control system which automatically locks and unlocks doors of a vehicle when the user carrying the portable device gets into and out of the vehicle, a plurality of portable devices are prepared and separately used in some cases when the security target such as a vehicle is used by a plurality of users.\nWhen the plural portable devices are used, however, the possibility that any of the portable devices is left inside the vehicle increases. If any of the plural portable devices is left inside the vehicle, the following problem occurs according to the conventional system.\nSince the left portable device is located inside the vehicle, the main device determines that the user as the owner of the portable device is still staying inside the vehicle and does not automatically set the security condition for the vehicle in some cases. When this occurs, the user does not notice that the portable device has been left inside the vehicle and leaves the vehicle in the non-security condition.\nIn another case, each of the plural users has his/her portable device and the security mechanism control system automatically locks the vehicle when any of the plural users gets off the vehicle. In still another case, the user locks the vehicle by directly giving a locking command to the vehicle.\nIn either of these cases, the vehicle is automatically unlocked by the approach of a person other than the users when any of the portable devices is left inside the vehicle. This is because the requirements that a person is approaching the vehicle and that the portable device is positioned near the main device are satisfied, and therefore the locking condition set for the vehicle is automatically cancelled.\nMoreover, when the plural portable devices are used in the system, such a case can happen where a user carrying the portable device moves away from the main device but another user stays near the main body.\nIn this case, the user leaving the vehicle considers that the vehicle has been automatically locked, but in reality the vehicle is not locked since the other user stays near the main device. Thus, the operation of the system depends on the behaviors of the users. There is a possibility that the requirement for unlocking is satisfied while the requirement for locking is satisfied. The above case is a result of the operation selected by the main body which has determined that the vehicle should be unlocked.\nTherefore, when the plural portable devices are used, the main device does not always set the condition as the user desires.\nAccording to the structures shown in Patent Reference Nos. 1 through 3, when the portable device left near the main body or carried by the user is positioned near the main body, the problems mentioned above cannot be solved. This is because the system does not have means for determining whether the portable device near the main body is the device left behind or the portable device carried by the user, nor means for appropriately notifying the user of the result of the determination.\nIn the structures shown in Patent Reference Nos. 1 through 3, therefore, in case of the plural portable devices are used, the vehicle cannot be locked or the locked condition is cancelled against the user's intension while the user is not noticing this situation."}
{"text": "Polynucleotide sequencing is the process of determining the nucleotide order of a given DNA or RNA fragment. These methods were first used in the sequencing of fragments, then individual genes, and have now been applied to whole genomes, from organisms as diverse as humans and E. coli. The desire for new sequence information has not been quenched by elucidation of these genomic sequences, however, as it has become apparent that individual variations in gene sequences play important roles in many physiological processes.\nAt the present time, most large scale sequencing efforts use the “chain termination”, or Sanger method. In this method, polymerization of DNA is initiated at a specific site on a template by using a short oligonucleotide primer. The primer is extended with DNA polymerase in the presence of four deoxynucleotides, along with a low concentration of a chain terminating nucleotide, for example a dideoxynucleotide. Limited incorporation of the chain terminating nucleotide results in a set of overlapping DNA fragments that are all terminated at a position corresponding to the dideoxynucleotide. The fragments are then size-separated by electrophoresis, e.g. by PAGE, capillary electrophoresis, etc. Detection of the fragments can utilize radioactive or fluorescent tags, and variations of the method use dye-labeled primers, dye-labeled terminator nucleotides, and the like.\nOther methods are known in the art that utilize size fractionation of fragments, for example the Maxam-Gilbert method, in which chemical reactions that selectively cleave DNA are used to generate the set of fragments. In addition, bulk techniques have been proposed that do not utilize sizing of fragments, including sequencing by hybridization, in which an array of short sequences of nucleotide probes is brought in contact with the target DNA sequence. A biochemical method determines the subset of probes that bind to the target sequence (the spectrum of the sequence), and a combinatorial method is used to reconstruct the DNA sequence from the spectrum.\nThere have been many proposals to develop new sequencing technologies based on single molecule measurements, for example by observing the interaction of particular proteins with DNA or by using ultra high-resolution scanned probe microscopy. Unlike conventional technology, their speed and read length would not be limited by the resolving power of electrophoretic separation. Single molecule sensitivity might permit direct sequencing of mRNA from rare cell populations or individual cells. A major obstacle has been the high data density of DNA. Scanned probe microscopes have not yet been able to demonstrate simultaneously the resolution and chemical specificity needed to resolve individual bases.\nBraslavsky et al. (2003) PNAS 100:3960-3964 imaged sequence information from a single DNA template as its complementary strand was synthesized. DNA template oligonucleotides were hybridized to a fluorescently labeled primer and bound to a solid surface via streptavidin and biotin with a surface density low enough to resolve single molecules. The primed templates were detected through their fluorescent tags, their locations were recorded for future reference, and the tags were photobleached. A combination of evanescent wave microscopy and single-pair fluorescence resonance energy transfer was used to reduce background noise. Labeled nucleotide triphosphates and DNA polymerase enzyme were then washed in and out of the flow cell while the known locations of the DNA templates were monitored for the appearance of fluorescence. It was shown that DNA polymerase was active on surface-immobilized DNA templates and can incorporate nucleotides with high fidelity, although only a few nucleotide residues were actually sequenced. A disadvantage of this method is that it used information from many single molecule images to generate an ensemble measurement from which sequence information was extracted. Therefore information from many molecules was averaged together to obtain a single measurement.\nAn alternative method has been proposed (Akeson et al. (1999) Biophys J. 77(6):3227-33) in which single molecules of a polynucleotide are driven through a nanopore channel by an applied electric field. During translocation, nucleotides within the polynucleotide are reported to pass through the pore in sequential, single-file order because of the limiting diameter of the pore. It has been suggested that this passage could provide a means for discrimination between pyrimidine and purine segments. At this time however, nanopore sequencing is a theoretical method, with only limited lab bench results.\nMethods of obtaining nucleotide sequences from small samples, down to single molecules, are of great interest for research, environmental and clinical purposes. The ability of such a method to provide information from a single molecule in the absence of averaging is desirable. The present invention addresses this issue.\nPublications\nAbbondanzieri et al. (2005) Nature 438:460-465 demonstrate direct observation of base pair stepping by RNA polymerase. Greenleaf et al. (2005) Phys. Rev. Lett. 95:208102 describe a passive all-optical force clamp for high-resolution laser trapping."}
{"text": "As the resin particles which are dispersed in the nonaqueous dispersion of resin particles, desired are those having a narrow particle size distribution and high dispersion stability. Above all, in such fields as the liquid developer for electrophotography, the liquid developer for electrostatic recording, the oil-based ink for inkjet printer, the ink for electronic paper, and the like, there are demands for resin particles having a small particle diameter and a narrow particle size distribution.\nIn addition, in the field of electrophotography, there is a strong desire for decreasing power consumption by printing machines such as printers and the like by fixing toners at lower temperatures.\nTherefore, in order to solve these problems, there has been proposed a method wherein a fatty acid monoester is added in the nonaqueous dispersion medium and, at the same time, a polyester resin is selected as the main component of the toner (see Patent Document 1). This method was intended to improve the fixing property of the toner by plasticizing the polyester resin with the fatty acid monoester. However, the toner particles obtained showed poor dispersion stability and were not entirely satisfactory in either particle size distribution or fixing property."}
{"text": "Insulation for spacecraft components helps to guard against the extreme conditions found in outer space, and to control thermal environments, for reliable operation of the spacecraft over long-durations. For example, because the exterior temperature of the spacecraft can vary by several hundred degrees depending on the spacecraft's exposure to solar radiation, such insulation may be used to thermally isolate the interior of the spacecraft or specific components to minimize thermal cycling. Multi-layer insulation (MLI) is considered the standard means for providing a thermal barrier for spacecraft components, and has been used in space-related applications for more than 50 years. MLI is typically composed of multiple flat layers of metallized Mylar or Kapton film, with a thin netting of an insulating polymer material, such as Dacron or Nomex, placed in between each film layer to minimize contact, and thus reduce thermal conduction between layers. However, the performance of MLI is strongly dependent on the manner in which it is formed and thereafter attached to the spacecraft, because areas where the MLI overlaps to produce folds, or areas having seams or penetrations, may dramatically reduce insulative performance. Because MLI's performance is so dependent on how well the individual flat panels fit together to form a covering that is free from compressive forces, the MLI blanket panels must be designed and sewn together in a custom manner, which is a time-consuming and expensive process. Moreover, MLI's insulative performance is difficult to predict due to its dependence on fit and the negative impact of seams or folds, which introduce unwanted conductive heat paths. Increasing the number of layers is a typical technique to compensate for this unpredictability."}
{"text": "Corn (Zea mays L.) is the most important and abundant crop produced in the United States. Corn is used as human food, livestock feed, and as raw material in industry. The food uses of corn include kernels for human consumption, dry milling products such as grits, meal and flour, and wet milling products such as corn starch, corn syrups, and dextrose. Corn oil recovered from corn germ is a by-product of both dry and wet milling industries. Both grain and non-grain portions of corn plants are used extensively as livestock feed, primarily for beef cattle, dairy cattle, hogs, and poultry.\nCorn is used to produce ethanol while corn starch and flour are used in the paper and textile industries. Corn is also used in adhesives, building materials, foundry binders, laundry starches, explosives, oil-well muds, and other mining applications. Plant parts other than the grain of corn are also used in industry; for example, stalks and husks are made into paper and wallboard and cobs are used for fuel and to make charcoal.\nThe goal of a corn breeder is to improve a corn plant's performance and therefore, its economic value by combining various desirable traits into a single plant. Improved performance is manifested in many ways. Higher yields of corn plants contribute to a more abundant food supply, a more profitable agriculture and a lower cost of food products for the consumer. Improved quality makes corn kernels more nutritious. Improved plant health increases the yield and quality of the plant and reduces the need for application of protective chemicals. Adapting corn plants to a wider range of production areas achieves improved yield and vegetative growth Improved plant uniformity enhances the farmer's ability to mechanically harvest corn.\nCorn is a monoecious plant, i.e., corn has imperfect flowers—male, pollen-producing flowers and separate female, pollen receiving flowers on the same plant. The male flowers are located at the top of the plant in the tassel, and the female flowers are located about midway up the stalk in the ear shoot. Each male flower has three anthers and each female flower includes a husk that envelops the cob and silks that emerge from the end of the cob and husks. Pollination is consummated by transfer of pollen from the tassels of the male flower to the silks of the female flowers.\nBecause corn has separate male and female flowers, corn breeding techniques take advantage of the plant's ability to be bred by both self-pollination and cross-pollination. Self-pollination occurs when pollen from the male flower is transferred to a female flower on the same plant. Cross-pollination occurs when pollen from the male flower is transferred to a female flower on a different plant.\nA plant is sib-pollinated (a type of cross-pollination) when individuals within the same family or line are used for pollination (i.e. pollen from a family member plant is transferred to the silks of another family member plant). Self-pollination and sib-pollination techniques are traditional forms of inbreeding used to develop new inbred corn lines but other techniques exist to accomplish inbreeding. New inbred corn lines are developed by inbreeding heterozygous plants and practicing selection for superior plants for several generations until substantially homozygous plants are obtained. During the inbreeding process with corn, the vigor of the lines decreases and after a sufficient amount of inbreeding, additional inbreeding merely serves to increase seed of the developed inbred. Inbred corn lines are typically developed for use in the production of hybrid corn lines.\nNatural, or open pollination, occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot and may include both self- and cross-pollination. Vigor is restored when two different inbred lines are cross-pollinated to produce the first generation (F1) progeny. A cross between two defined homozygous inbred corn plants always produces a uniform population of heterozygous hybrid corn plants and such hybrid corn plants are capable of being generated indefinitely from the corresponding inbred seed supply.\nWhen two different, unrelated inbred corn parent plants are crossed to produce an F1 hybrid, one inbred parent is designated as the male, or pollen parent, and the other inbred parent is designated as the female, or seed parent. Because corn plants are monoecious, hybrid seed production requires elimination or inactivation of pollen produced by the female parent to render the female parent plant male sterile. This serves to prevent the inbred corn plant designated as the female from self-pollinating. Different options exist for controlling male fertility in corn plants such as manual or mechanical emasculation (or detasseling), genetic male sterility, and application of gametocides. Incomplete removal or inactivation of the pollen in the female parent plant provides the potential for inbreeding which results in the unwanted production of self-pollinated or sib-pollinated seed. Typically, this seed is unintentionally harvested and packaged with hybrid seed.\nThe development of new inbred corn plants and hybrid corn plants is a slow, costly interrelated process that requires the expertise of breeders and many other specialists. The development of new hybrid corn varieties in a corn plant breeding program involves numerous steps, including: (1) selection of parent corn plants (germplasm) for initial breeding crosses; (2) inbreeding of the selected plants from the breeding crosses for several generations to produce a series of inbred lines, which individually breed true and are highly uniform; and, (3) crossing a selected inbred line with an unrelated line to produce the F1 hybrid progeny having restored vigor.\nInbred corn plants and other sources of corn germplasm are the foundation material for all corn breeding programs. Despite the existence and availability of numerous inbred corn lines and other source germplasm, a continuing need still exists for the development of improved germplasm because existing inbred parent corn lines lose their commercial competitiveness over time. To protect and to enhance yield production, trait technologies and seed treatment options provide additional crop plan flexibility and cost effective control against insects, weeds and diseases, thereby further enhancing the potential of hybrids with UWC04 as a parent."}
{"text": "Field of the Invention\nThe present application relates generally to electrical machines including, but not limited to electrical motors and electrical generators and more particularly, but not by way of limitation to electrical machines that are co-axially arranged with a magnetic gear.\nHistory of the Related Art\nElectrical machines have been utilized for a variety of purposes. Typically, the cost of an electric machine is proportional to the torque with which it must interact. For applications requiring a large torque at a low speed there are two generally accepted options. First, a high speed, low torque machine could be utilized in conjunction with a mechanical gear. Second, a larger low speed, high torque direct-drive machine could be utilized. Mechanical gears are generally viewed as unattractive options as they require considerable maintenance, produce acoustic noise, and have a shorter useful life.\nMagnetic gears have attracted considerable attention as a possible replacement for traditional, mechanical gears. Unlike mechanical gears, which rely on the physical interaction between teeth, magnetic gears create a gearing action through the modulated interaction between magnetic flux generated by two means of inducing magnetic flux with different pole pair counts. Magnetic gears exhibit generally contactless operation and, as such, facilitate reduced maintenance, improved reliability, decreased acoustic noise, and physical isolation between input and output shafts."}
{"text": "1. Field of the Invention\nThe present invention relates to a holding apparatus; in particular, the present invention relates to a holding apparatus for a projector that can adjust the position of a laser source according to the disposing surface.\n2. Description of the Prior Art\nProjection device has the advantage of outputting large size image. With the technical development of the image quality, projectors have been widely used in conference rooms, exhibitions, or classrooms. There are some products combining the projector with a laser touch module for users to write directly on the projection image. Specifically, the technology uses a laser touch module generating laser beams to form a laser screen which covers the projection image. When an object touches the projection image, the laser beams will be reflected. The lens of the projector will identify the relative position to achieve the direct writing operation.\nHowever, the present products combining projectors with the laser touch module require the laser touch module directly installed on a wall or at an inner side of a border of a white board. After deducting safety distance of the laser touch module, the laser screen cannot match the projection image. In other words, writing actions cannot be sensed at some locations on the projection image resulted from the configuration of the traditional laser touch module. When users write on the projection image, the writing actions may not be identified. As such, the present products combining projectors with the laser touch module need to be improved."}
{"text": "In a conventional direct-mode communication system, apparatuses can communicate with each other directly, so messages can be exchanged more simply and rapidly. However, usually no complete security protocol is available for such direct-mode communications, so the direct-mode communications are liable to unauthorized use by malicious persons to forge identifications and to steal information. Accordingly, when the discovery interactive technology is used for the conventional direct-mode communication apparatuses, the security and the stability of these apparatuses will become a significant potential trouble.\nAdditionally, if it is desired to transmit interactive information to a user apparatus when the discovery interactive technology is used in the conventional direct-mode communication system, usually an interactive information request is firstly transmitted by the user apparatus to a nearby serving apparatus. Then, the serving apparatus transmits the interactive information request to a back-end core network so that the back-end core network broadcasts interactive information to the user apparatus. However, information transmission in this manner is not only relatively complicated, but is also likely to cause a burden on network resources due to redundant information transmission.\nAccordingly, an urgent need exists in the art to enhance the security of the direct-mode communication system and increase the convenience of message propagation."}
{"text": "1. Field of the Invention\nThe present invention relates to a recording apparatus and, more specifically, an audiovisual recording apparatus.\n2. Description of the Related Art\nHard disk drive (HDD) recorders configured to receive television broadcasts and to record them onto HDDs have been available. Recently, video cameras equipped with HDDs have also been available.\nIn such HDD recorders or HDD-equipped video cameras, HDDs are used to temporarily record audio-video data such as audio-video data of recorded television broadcasts or captured audio-video data. Audio-video data temporarily recorded on HDDs is usually duplicated and recorded onto digital versatile disks (DVDs) by users using personal computers (PCs), built-in DVD drives, or the like.\nFor example, U.S. Patent Laid-Open No. 2004/0246534 discloses a video camera capable of duplicating audio-video data to a DVD by a single touch operation in cooperation with an application operating on a PC.\nWhen recording audio-video data recorded in the manner described above onto external media, users search on PCs for a desired audio-video data file from a large number of audio-video data files and enter an instruction to record it."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of oscillating an ultrasonic vibrator for use in ultrasonically cleaning (including deburring) workpieces immersed in a cleaning solution.\n2. Description of the Prior Art\nFor ultrasonically cleaning workpieces immersed in a cleaning solution in a cleaning tank, it has been customary to apply a periodic voltage signal to an ultrasonic vibrator having a piezoelectric element, the periodic voltage signal having a frequency equal to the natural frequency of the ultrasonic vibrator, to oscillate the ultrasonic vibrator at its natural frequency for thereby radiating an ultrasonic energy into the cleaning solution. The radiated ultrasonic energy produces a cavitation in the cleaning solution, which generates shock waves to clean and deburr the workpieces immersed in the cleaning solution.\nIt is generally known that the cavitation in the cleaning solution appears at a depth depending on the frequency of the radiated ultrasonic energy, i.e., the natural frequency (resonant frequency) of the piezoelectric element of the ultrasonic vibrator. More specifically, when the ultrasonic energy is radiated from the bottom of the cleaning tank toward the surface level of the cleaning solution in the cleaning tank, the cavitation is produced intensively at a depth equal to a quarter wavelength, and also at depths positioned successively at quarter wavelength intervals from that depth toward the bottom of the cleaning tank.\nFor uniformly cleaning and deburring the workpieces immersed in the cleaning solution, it is preferable to generate the cavitation uniformly in the cleaning solution without being dispersed in the cleaning solution. To generate the cavitation uniformly in the cleaning solution, it is desirable to radiate the ultrasonic energy at a higher frequency. It is also generally known that the higher the frequency of the radiated ultrasonic energy, the more the ultrasonic energy is attenuated in the cleaning solution, resulting in a lowered cavitation effect. For effective cleaning or deburring of the workpieces, therefore, it is preferable to radiate the ultrasonic energy at a lower frequency. Since the generation and effect of the cavitation vary depending on the frequency of the ultrasonic energy, the frequency of the ultrasonic energy should be selected in view of the purpose for which the workpieces are to be cleaned and the degree to which the workpieces are to be cleaned. For example, if a stronger cleaning capability is desirable, then the ultrasonic energy should be applied at a lower frequency. If the workpieces to be cleaned are fragile, then the ultrasonic energy should be applied at a higher frequency in order to prevent the workpieces from being damaged by the cavitation.\nHowever, where an ultrasonic vibrator having a single natural frequency is oscillated at the natural frequency, the above requirements cannot be satisfied under various conditions.\nOne solution has been to employ an ultrasonic vibrator having a plurality of piezoelectric elements having respective different natural frequencies, and repeatedly apply a plurality of signals having frequencies equal to the natural frequencies to the respective piezoelectric elements for respective periods of time. Therefore, ultrasonic energies are radiated at different frequencies from the single ultrasonic vibrator into the ultrasonic solution.\nWhen the ultrasonic energies are radiated into the ultrasonic solution, cavitations are produced at relatively close depths, respectively, in the cleaning solution. As a result, the cavitations are distributed comparatively uniformly in the cleaning solution, and it is possible to obtain an effective cavitation effect primarily based on those ultrasonic energies which have lower frequencies. A suitable choice of periods of time for which the ultrasonic energies having different frequencies are radiated is effective to serve different purposes for which workpieces are to be cleaned.\nThe ultrasonic vibrator with plural piezoelectric elements having respective different natural frequencies, however, is difficult and expensive to manufacture. Another problem is that the cavitation distribution becomes unstable because the natural frequencies of the piezoelectric elements tend to vary due to the heat produced thereby when the ultrasonic vibrator is oscillated. Consequently, it has been difficult to clean and deburr the workpieces uniformly with the cavitations."}
{"text": "Light fixtures, or luminaires, are used with electric light sources to provide an aesthetic and functional housing in both interior and exterior lighting applications. For example, high bay luminaires can be used in larger open indoor environments such as heavy industrial settings, warehouses, gyms, churches, and shopping malls. Conventional high bay lighting fixtures for commercial and industrial applications are often mounted or suspended from ceiling joists high above the floor.\nRecently, lighting fixtures have begun using light emitting diodes (LEDs) as a light source. The use of LEDs comes with unique light distribution and thermal management requirements for both the light emitters and the control components required to run the light fixtures. These considerations can lead to complex housing and heat dissipation designs in an attempt to balance performance and aesthetic characteristics."}
{"text": "This invention relates to an adsorbent material for recovering and removing objectionable substances, in particular, arsenic that are contained in environmental water and liquid wastes such as waste water from plants. One or more embodiments of the invention also relate to a method of synthesizing the adsorbent material.\nResearch and development efforts have recently been made on materials capable of trapping metals contained in environmental water bodies such as rivers and the sea and this has led to the discovery that cation exchange resins using phosphoric acid as exchange groups can adsorb metal ions present in rivers, lakes and wastewater from plants (see, for example, Akinori Joh et al., “Cation exchange resins using phosphoric acid as exchange groups—Their selectivity for metal ions and applications” in PHOSPHORUS LETTER, Japanese Association of Inorganic Phosphorus Chemistry, Feb. 1, 2001, vol. 40, pp. 16-21).\nFurther, the present inventors developed a metal adsorbent material that had a monomer with phosphoryl groups grafted onto a polymeric substrate (see, for example, Japanese Patent Application No. 2002-262502).\nIn the art of recovering and removing objectionable substances in the environment, particularly arsenic, two major methods have so far been practiced, one relying upon coagulating sedimentation and the other using chelating resins (see, for example, Xiaoping Zhu et al., “Removal of arsenic(V) by zirconium(IV)-loaded phosphoric acid chelating resin” in “SEPARATION SCIENCE AND TECHNOLOGY”, America, Marcel Dekker, Inc., 2001, 36(14), pp. 3175-3189).\nHowever, the conventional adsorbent materials can adsorb arsenic only slowly. If arsenic is recovered and removed by coagulating sedimentation or with the aid of adsorbent materials in bead form, considerable inconvenience in handling has been met during the process of removal or in subsequent operations.\nIn addition, the conventional arsenic adsorbent materials are mostly synthesized by common radical polymerization and their structure for adsorption of arsenic is so unstable that it is prone to leak out even if it is adsorbed.\nAccording to the Basic Environment Law which provides for the water quality guidelines for public waters, arsenic should not be discharged at concentrations higher than 0.1 ppm and its content in the environment should not exceed 0.01 ppm. Thus, the removal of arsenic is absolutely necessary."}
{"text": "Systems and processes disclosed herein relate to the recovery of sorbates that are removed from process streams by adsorption.\nThe separation of specific components of process streams in certain applications can be accomplished by the use of adsorption to remove the desired component, followed by recovery of the component during regeneration of the adsorbent. Process streams are typically liquid or gaseous, and can have a wide variety of compositional components depending upon the industrial application of the process.\nPacked beds of adsorbent materials are typically used in adsorption processes. Adsorbent materials are generally in the form of spherical beads, or pellets. Adsorbent materials are typically oxygen-containing compounds, carbon-containing compounds, or polymer-based compounds. Oxygen-containing compounds can be, for example, hydrophilic and polar, including materials such as silica gel and zeolites. Carbon-based compounds can be, for example, hydrophobic and non-polar, including materials such as activated carbon and graphite. Polymer-based compounds can be, for example, polar or non-polar functional groups in a porous polymer matrix.\nTypical adsorption processes utilizing packed beds can be thermal (temperature) swing adsorption (TSA) processes or pressure swing adsorption (PSA) processes. In operation, a process stream is introduced into a packed bed, and the adsorbent material contained therein removes a desired component, known as the sorbate, from the stream as it filters through the packed bed. After a given time period, the adsorbent material becomes saturated with the sorbate, and the adsorption process must be halted in order to regenerate the adsorbent and remove the sorbate. PSA processes utilize a de-pressurized regeneration gas that is introduced to the packed bed in a direction reverse to the flow of the process stream. After a regeneration cycle is complete, a new adsorption cycle can begin. TSA processes utilize heat to remove the sorbate from the adsorbent material. The heat in a typical TSA process is added through the regeneration stream. Thus, the regeneration or purge gas can be utilized to supply all of the heat required to heat the vessel, bed supports, and the adsorbent, as well a providing the energy to desorb the sorbate from the adsorbent. A large volume of regeneration gas is necessary to accomplish this task, which results in a very dilute regeneration effluent stream, which increases the cost of processes utilized to recover the sorbate from the regeneration effluent stream."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a deep trench capacitor, and more specifically to the process of fabricating and the device of a trench capacitor, providing increased capacitance.\n2. Description of the Prior Art\nAn important design goal in the design and fabrication of deep trench capacitors is reduced size. Another design goal in designing and fabricating deep trench capacitors is increased capacitance.\nIn accordance with some conventional processes for fabricating deep trench capacitors, capacitors are formed above the substrate. This type of structure utilizes some of the limited space above the substrate. Therefore, the available space for fabricating other devices is reduced. In order to provide space for fabricating other devices, and in order to increase the density of devices, the size of the trench capacitor must be reduced or the capacitor may be built within the substrate. A trench type capacitor is a capacitor that is fabricated within the substrate.\nFIG. 1 shows a transverse cross-sectional view of a typical prior art trench type capacitor at 10. A trench 12 having sidewalls and a bottom is formed within a semiconductor substrate 14. The trench 12 is typically formed in the substrate 14 by a process of patterning and etching the substrate 14. Subsequently, a first conductor 18 is formed in a portion of the substrate 14 that envelopes a portion of the trench 12 by doping the inside walls and bottom surface of the trench 12 with arsenic ion (or phosphoric ion). The first conductor 18 is used as a bottom conductor, or xe2x80x9cbottom electrode platexe2x80x9d, of the conventional trench type trench capacitor 10. After forming the first conductor 18 by doping the semiconductor substrate 14 as described above, a dielectric layer 22 is formed typically by depositing a silicon nitridexe2x80x94silicon oxide layer (commonly referred to as an NO layer). After forming the dielectric layer 22, a second conductor 26 is formed by depositing amorphous silicon into the trench 12 and over the dielectric layer 22. The second conductor 26 provides a top conductor, or xe2x80x9ctop electrode platexe2x80x9d, of the conventional trench type capacitor 10.\nOne problem associated with the depicted prior art deep trench capacitor 10 is that it provides only a limited amount of capacitance. The capacitance of the deep trench capacitor 10 is proportional to the surface areas of the first and second plates 18 and 26. It is possible to increase the capacitance of the deep trench capacitor 10 by increasing the size of the plates. However, increasing the size of the plates also increases the total size of the trench type capacitor which defeats the design goal of minimizing the size of the trench capacitor.\nIt is an object of the present invention to provide a deep trench capacitor having increased capacitance, and requiring a decreased amount of space. What is needed is a design and method of fabricating a deep trench capacitor having increased capacitance, and decreased size.\nIt is an object of the present invention to provide a deep trench capacitor having increased capacitance, and decreased size.\nIt is another object of the present invention to provide a process of fabricating a deep trench capacitor having increased capacitance, and requiring a decreased amount of space.\nBriefly, a presently preferred embodiment of the present invention provides a process of fabricating a trench capacitor. The process includes the steps of: depositing a nitride masking layer over a top surface of a semiconductor substrate; patterning and etching to remove a portion of the nitride masking layer and a portion of the substrate to form an exposed deep trench in the substrate, and to expose an edge of the nitride masking layer, the trench having sidewalls and a bottom surface; forming an oxide fill plug to fill a bottom portion of the trench, the oxide fill plug extending from the bottom surface of the trench to a predetermined depth level below the top surface of the substrate; etching a portion of the sidewalls to form an etched back portion of the sidewalls located between the oxide fill plug and the top surface of the substrate; forming a collar insulating layer within the etched back portion of the sidewalls; forming a protection layer over the exposed edge of the nitride masking layer and over the collar insulating layer; and removing the oxide fill plug from the trench.\nThe process also includes the steps of: doping a region of the substrate enveloping the bottom portion of the trench; depositing a spacer insulating layer over the protection layer, the sidewalls, and the bottom surface of the trench; removing a portion of the spacer insulating layer to expose a substantially central portion of the bottom surface of the trench; depositing a conducting layer over the nitride masking layer, the spacer insulating layer, and the exposed central portion of the bottom surface, the conducting layer and the doped region of the substrate being in electrical contact and forming a first plate of the capacitor; depositing a photo resist fill plug into the trench and over the conducting layer, the photo resist fill plug providing protection of a bottom portion of the conducting layer, and leaving an unprotected portion of the conducting layer; removing the unprotected portion of the conducting layer to expose a top surface of the conducting layer; and removing the photo resist fill plug to expose inner walls of the conducting layer.\nThe process further includes the steps of: removing the spacer insulating layer to expose outer walls of the conducting layer, and leaving empty space between the outer walls and the sidewalls; removing the protection layer; forming a dielectric layer over the sidewalls of the trench and over the top surface, the inner walls, and the outer walls of the conducting layer; and depositing a conducting plug over the dielectric layer and filling the trench, the conducting plug forming at least a portion of a second plate of the capacitor.\nAn important advantage of a deep trench capacitor formed in accordance with the present invention is that it provides increased capacitance, and requires a decreased amount of space thereby optimizing the use of silicon real estate."}
{"text": "Field of the Invention\nThe present invention relates to an automated assembly apparatus and an automated assembly method, which use an assembly robot.\nDescription of the Related Art\nHitherto, in an automated assembly apparatus which uses a robot, a Z-axis of the robot gripping an assembly component moves up and down to assemble the assembly component to a component to be assembled.\nAutomated assembly apparatuses have been required to achieve reduced assembling operation time. Japanese Patent Application Laid-Open No. H5-235598 describes a component mounting apparatus which suctions a component by a suction nozzle and moves horizontally with respect to a substrate to mount the component at a predetermined position on the substrate.\nThe component mounting apparatus described in Japanese Patent Application Laid-Open No. H5-235598 is capable of reducing the component mounting time to a certain extent. However, the Z-axis has to be raised to a height which clears the maximum height of an interfering object on the movement path when XY-axes move. In addition, when mounting a component, the Z-axis has to be lowered by the same height it has been lifted to clear the interfering object. This has been posing a problem in that the time for mounting a component increases as the height of an interfering object on the movement path increases."}
{"text": "The invention relates to a lubricator for underbalanced drilling.\nThere are two techniques that typically are used to drill a borehole in a formation: an overbalanced drilling technique and underbalanced drilling technique. In overbalanced drilling, fluid in an annulus of a well is used to exert a pressure that is greater than the formation pressure. Thus, the pressure that is exerted by the annulus fluid keeps formation fluids from exiting the well. A drawback to this technique is that mud particles typically are added to the annulus fluid to increase its weight (and thus, increase its downhole pressure), and these mud particles tend to clog up openings in the formation. Thus, the formation may be damaged by overbalanced drilling, and after drilling, cleanup of the well may be needed before production begins. The well may also need to be tested after overbalanced drilling to check for formation damage.\nUnlike overbalanced drilling, underbalanced drilling typically does not damage the formation damage and typically maximizes reservoir inflow. In underbalanced drilling, heavy annulus fluid is not used to suppress the formation pressure. Instead, a blowout preventer, or snubbing unit, is used to seal off the drill string at the surface of the well. However, this arrangement may also present difficulties. For example, when drilling at shallow depths or retrieving the drill string, the upward force from the formation pressure may exceed the weight of the drill string and thus, may force the drill string out of the borehole. As a result, retrieving the drill string may consume a considerable amount of time and present a significant danger.\nThus, there is a continuing need for an arrangement to address one or more of the problems that are stated above.\nIn an embodiment of the invention, a system usable with a subterranean well includes a tubing and a lubricator. The tubing is adapted to receive a drill string in a passageway of the tubing, and the lubricator is located downhole and is connected to the tubing. The lubricator is adapted to be remotely operable from a surface of the well to control fluid communication between the passageway located above the lubricator and a formation located beneath the lubricator.\nIn another embodiment of the invention, an apparatus that is usable with a downhole tool that has a passageway includes a sleeve and a controller. The controller selectively moves the sleeve into the passageway to protect a portion of the downhole tool from a downhole fluid.\nAdvantages and other features of the invention will become apparent from the following description, drawing and claims."}
{"text": "Many substituted pyrido(3,2-d)pyrimidine derivatives are known in the art. For instance, pyrido(3,2-d)pyrimidine derivatives with various substituents on positions 2, 4 and 6 (using the standard atom numbering for the pyrido(3,2-d)pyrimidine moiety) are known with biological activities such as competitive inhibition of pteroylglutamic acid, inhibition of thrombocyte aggregation and adhesiveness, antineoplastic activity, inhibition of dihydrofolate reductase and thymidylate synthase, e.g., from U.S. Pat. Nos.2,924,599, 3,939,268, 4,460,591, 5,167,963 and 5,508,281.\nPyrido(3,2-d)pyrimidine derivatives with various substituents on positions 2, 4, 6 and 7 (using the standard atom numbering for the pyrido(3,2-d)pyrimidine moiety), some of them with biological activities, are also known e.g., from U.S. Pat. No. 5,521,190, U.S. patent application publication No. 2002/0049207, U.S. patent application publication No. 2003/0186987, U.S. patent application publication No. 2003/0199526, U.S. patent application publication No. 2004/0039000, U.S. patent application publication No. 2004/0106616, U.S. Pat. Nos. 6,713,484, 6,730,682 and 6,723,726.\nU.S. Pat. No. 5,654,307 discloses pyrido(3,2-d)pyrimidine derivatives substituted on position 4 with monoarylamino or monobenzylamino, and on positions 6 and 7 with substituents each independently selected from the group consisting of lower alkyl, amino, lower alkoxy, mono- or dialkylamino, halogen and hydroxy. WO 01/083456 discloses pyrido(3,2-d)pyrimidine derivatives substituted on position 4 with morpholinyl and on position 2 with hydroxyphenyl or morpholinoethoxyphenyl, having PI3K and cancer inhibiting activity. U.S. Pat. No. 6,476,031 discloses substituted quinazoline derivatives, including (in reaction scheme 5) a series of pyrido(3,2-d)pyrimidine derivatives which are substituted on position 4 with hydroxy, chloro or an aryl, heteroaryl (including pyridyl, pyrimidyl, indolyl, benzimidazolyl, benzotriazolyl, isoquinolyl, quinolyl, benzothiazolyl, benzofuranyl, thienyl, furyl, pyrrolyl, thiazolyl, oxazolyl, imidazolyl), cycloaliphatic or cycloheteroaliphatic group being optionally spaced from the pyrido(3,2-d)pyrimidine ring by a linker such as NH.\nWO 02/22602 and WO 02/22607 disclose pyrazole and triazole compounds, including 2-(1-trifluoromethylphenyl)-4-fluorobenzopyrazolyl-pyrido(3,2-d)pyrimidine and 2-(1-tri-fluoromethylphenyl)-4-methyltriazolyl-pyrido(3,2-d)pyrimidine being useful as protein kinase inhibitors. WO 03/062209 discloses pyrido(3,2-d)pyrimidine derivatives substituted on position 7 with aryl or heteroaryl and on position 4 with monoarylamino or monoheteroarylamino and which may further be substituted on positions 2 and/or 6, being useful as capsaicin receptor modulators.\nWO 2006/069805 discloses pyrido(3,2-d)pyrimidine derivatives substituted on position 6 with aryl or heteroaryl and on both positions 2 and 4 with monoalkylamino, monocycloalkylamino, monoarylamino or monoarylalkylamino, and which may further be substituted on position 7, being useful in the treatment of a disease mediated by phosphodiesterase-4 activity. WO 2006/135993 discloses di- and tri-substituted pyrido(3,2-d)pyrimidine derivatives being highly active and virus-specific anti-flaviridae agents which are especially useful in the treatment of hepatitis C.\nMore recently, novel pyrido(3,2-d)pyrimidine derivatives have been disclosed in WO2008/009078, WO2008/009076, WO2008/009079, WO2008/009077, WO2008/043149, PCT/EP2007/011494, PCT/EP2007/011495 and PCT/EP2007/011496.\nNotwithstanding the discoveries of new compounds disclosed in the above patent applications, there remains a continuous need in the art for specific and highly therapeutically active compounds for preventing or treating infections due to Flaviridae and pathologic conditions associated therewith, especially hepatitis C. In particular, there is a need in the art to provide drugs which are active against hepatitis C in order to replace existing drugs having significant side effects and to decrease treatment costs.\nHepatitis is an inflammation of the liver that is most often caused by infection with one of three viruses known as hepatitis A, B or C. HCV is a representative and highly significant member of the Flaviviridae family, a family of positive-strand RNA viruses. Current standard of care for HCV treatment is a combination of (pegylated) interferon alpha and the antiviral drug ribavirin for a period of 24 or 48 weeks, depending upon the viral genotype.\nBest results have been achieved with the combination of weekly subcutaneous injections of long-acting peginterferon alpha and oral ribavirin daily. Interferons are substances naturally released by cells in the body after viral invasion. Interferon alfa-2b and peginterferon alfa-2b are synthetic versions of these substances. The protein product is manufactured by recombinant DNA-technology. Second generation interferons are further derivatized by binding to inert polyethylene glycol, thereby altering the pharmacokinetic properties. Ribavirin is a nucleoside analogue, which disrupts viral replication of hepatitis C virus (HCV).\nThe most common side effects of HCV treatment with (pegylated) interferon include: a decrease in white blood cells and platelets, anemia, nausea, diarrhea, fever, chills, muscle and joint pain, difficulty in concentrating, thyroid dysfunction, hair loss, sleeplessness, irritability, mild to serious depression, and rarely, suicidal thoughts. Other serious adverse events include bone marrow toxicity, cardiovascular disorders, hypersensitivity, endocrine disorders, pulmonary disorders, colitis, pancreatitis, and ophthalmologic disorders (eye and vision problems). (Pegylated) interferon may also cause or make worse fatal or life-threatening neuropsychiatric, autoimmune, ischemic, and infectious disorders. Patients with persistently severe or worsening signs or symptoms of these conditions are advised to stop therapy.\nThe most common side effect of HCV treatment with ribavirin is anaemia, which can be treated with erythropoietin. Other side effects include mood swings, irritability, anxiety, insomnia, abdominal pain, nervousness, breathlessness, rash, hair loss, dry skin, nausea, diarrhoea, loss of appetite, dizziness and weight loss. Ribavirin can also cause birth defects. Ribavirin should not be taken in combination with certain HIV drugs such as, but not limited to, didanosine, since lactic acidosis with fatal hepatic steatosis (fatty liver) may occur. Special attention should be taken for treatment with HIV co-infection.\nThere is a strong need in the art to improve, or to provide alternatives to, the existing prophylactic or therapeutic solutions to infections by a virus of the Flaviridae family, more specifically HCV infection. In particular there is still a need in the art for providing alternative synthetic molecules having significant HCV replication inhibiting activity and improved druggable properties, particularly, there is a need to provide improved pyrido(3,2-d)pyrimidines. There is also a need in the art for providing effective inhibiting molecules which are free from the significant drawbacks of current drugs, like pegylated interferon and ribavirin. Meeting these various needs in the art constitutes a main goal of the present invention."}
{"text": "1. Field of the Invention\nThe invention relates to an electrically active particle having a virus with bound nanoparticles and molecular wires.\n2. Description of the Prior Art\nImprovements in the speed and power of modern computers are produced mostly through increasing the number of available transistors in an integrated circuit. Although Moore's Law, which states that the number of transistors that can be fabricated on a silicon integrated circuit will double every 18 to 24 months, has been valid for the past 40 years, eventually the feature size of the transistors will small enough that quantum effects will prevent their operation. One possible alternative to traditional silicon-based technology is to use functional molecules as the building blocks of electronic devices. Instead of carving smaller and smaller patterns into silicon wafers, this bottom-up approach uses different types of molecules with functions such as wires, switches, and diodes, to build electronic circuits, and thus increasing the theoretical device density by up to 6 orders of magnitude.\nAlthough a large number of molecules have been synthesized for their electronic properties, the problem remains in connecting these active molecules in such a way as to build logic circuits and functional electronic devices. Two of the approaches are (1) a crosswire device in which the electronically addressable molecules are sandwiched between the junction and (2) an array of containers with input and output pins on opposing sides with the molecules of interest bridging the two sets of pins."}
{"text": "1. Field of the Invention\nThe present invention relates to a sewing machine and, more particularly, to a sewing machine capable of automatically carrying out anti-loosening-stitching and thread-cutting.\n2. Description of the Related Art\nA known sewing machine capable of automatically carrying out back-stitching and thread-cutting after completing a stitch pattern is provided with an automatic back-stitching mode selector switch and an automatic thread-cutting mode selector switch respectively for setting the back-stitching activated mode and the thread-cutting activated mode. The automatic back-stitching mode selector switch and the automatic thread-cutting mode selector switch are operated to set the back-stitching activated mode and the thread-cutting activated mode and to cancel the same individually (i.e., to set the back-stitching deactivated mode and the thread-cutting deactivated mode). When the completion of a stitch pattern requires automatic back-stitching and automatic thread-cutting or when canceling the set back-stitching mode and the set thread-cutting mode, the automatic back-stitching mode selector switch and the automatic thread-cutting mode selector switch are operated accordingly to set the back-stitching activated mode and the thread-cutting activated mode or to set the back-stitching deactivated mode and the thread-cutting deactivated mode."}
{"text": "An aircraft may perform a variety of different missions. Some missions may subject the aircraft to more stress than other missions. For example, some missions may include high-stress maneuvers, whereas other missions may focus primarily on straight-and-level flight.\nOne example of an aircraft is a rotorcraft. A rotorcraft may include one or more rotor systems. One example of a rotorcraft rotor system is a main rotor system. A main rotor system may generate aerodynamic lift to support the weight of the rotorcraft in flight and thrust to counteract aerodynamic drag and move the rotorcraft in forward flight. Another example of a rotorcraft rotor system is a tail rotor system. A tail rotor system may generate thrust in the same direction as the main rotor system's rotation to counter the torque effect created by the main rotor system."}
{"text": "The present invention generally relates to fabrication of large composite structures and, more particularly, to automated composite lay-up of large aircraft fuselage sections.\nThe structural performance advantages of composites, such as carbon fiber epoxy and graphite bismaleimide (BMI) materials, are widely known in the aerospace industry. Aircraft designers have been attracted to composites because of their superior stiffness, strength, and radar absorbing capabilities, for example. As more advanced materials and a wider variety of material forms have become available, aerospace usage of composites has increased. Automated tape layer technology has developed to become a widely used automated process for fabrication of large composite structures such as wing panels and empennage. Current tape layer technology has been improved to offer flexibility in process capabilities required for a wide variety of aerospace components. As aerospace industry tape laying applications achieve material lay up rates, for example, that may help control the manufacturing cost of large composite structures, new and innovative applications for tape layers may be defined, such as the automated tape lay-up of large aircraft fuselage sections, for example, 15 to 20 feet in diameter.\nAutomated tape laying machines typically are gantry style machines that may have, for example, ten axes of movement with 5-axis movement on the gantry and 5-axis movement on the delivery head. A typical automated tape layer consists of a gantry structure (parallel rails), a cross-feed bar that moves on precision ground ways, a ram bar that raises and lowers the material delivery head, and the material delivery head which is attached to the lower end of the ram bar. Commercial tape layers are generally configured specifically for lay up of flat or mildly contoured laminate applications using either flat tape laying machines (FTLM) or contour tape laying machines (CTLM). On a gantry style tape layer, tooling (or a flat table) is commonly rolled under the gantry structure, secured to the floor, and the machine delivery head is then initialized to the lay up surface.\nFIG. 1 provides an illustration of a typical tape laying machine material delivery head 100. Delivery heads for FTLM and CTLM machines are basically the same configuration as that of delivery head 100 shown in FIG. 1. The delivery heads on commercial automated tape layers are typically configured to accept material widths of 75 mm (3″), 150 mm (6″), and 300 mm (12″). Flat tape layers typically use material in 150 mm (6″) and 300 mm (12″) widths. Contour tape layers typically use material in 75 mm (3″) and 150 mm (6″) widths. CTLM systems normally use the 3″ or 6″ wide material when laying up off flat plane contour surfaces. Material 102 for tape layers generally comes in large diameter spools. The tape material 102 has a backing paper 106, which must be extracted as the prepreg (resin pre-impregnated fiber) is applied to the tool surface 108. The spool of material typically is loaded into the delivery head supply reel 104 and threaded through the upper tape guide chute and past the cutters 110. The material 102 then passes through the lower tape guides, under the segmented compaction shoe 112, and onto a backing paper take up reel 114. The backing paper is extracted and wound on a take up roller of paper take up reel 114. The delivery head 100 makes contact with the tool surface 108 and the tape material 102 is “placed” onto the tool surface 108 with compaction pressure. The tape laying machine typically lays tape on the tool surface 108 in a computer programmed path (course), cuts the material 102 at a precise location and angle, lays out tail, lifts delivery head 100 off the tool surface 108, retracts to the course start position, and begins laying the next course. The delivery head 100 may have an optical tape flaw detection system that signals the machine control to stop laying tape material 102 when a flaw has been detected. The delivery head 100 also typically has a heating system 116 that heats the prepreg materials to increase tack levels for tape-to-tape adhesion. Heated tape temperatures generally range from 80 F to 110 F.\nFiber placement is a similar process in which individual prepreg fibers, called tows, are pulled off spools and fed through a fiber delivery system into a fiber placement head, which is similar to delivery head 100 shown in FIG. 1. In the fiber placement head, tows may be collimated into a single fiber band and laminated onto a work surface, which can be mounted between a headstock and tailstock. When starting a fiber band or course, the individual tows are fed through the head and compacted onto a surface—such as surface 108. As the course is being laid down, the head 100 can cut or restart any of the individual tows. This permits the width of the fiber band to be increased or decreased in increments equal to one tow width. Adjusting the width of the fiber band eliminates excessive gaps or overlaps between adjacent courses. At the end of the course, the remaining tows may be cut to match the shape of the ply boundary. The head may then be positioned to the beginning of the next course. During the placement of a course, each tow is dispensed at its own speed, allowing each tow to independently conform to the surface 108 of the part. Because of this, the fibers are not restricted to geodesic paths. They can be steered to meet specified design goals. A rolling compaction device, combined with heat for tack enhancement, laminates the tows onto the lay-up surface 108. This action of pressing tows onto the work surface (or a previously laid ply) adheres the tows to the lay-up surface 108 and removes trapped air, minimizing the need for vacuum debulking. It also allows the fiber to be laid onto concave surfaces.\nA fiber placement head, like the tape laying head, may be provided with several axes of motion, using an arm mechanism, for example, and may be computer numeric controlled. The axes of motion may be necessary to make sure the head 100 is normal to the surface 108 as the machine is laminating tows. The machine may also have a number of electronic fiber tensioners, which may be mounted, for example, in an air conditioned creel. These tensioners may provide individual tow payout and maintain a precise tension. The head 100 may precisely dispense, cut, clamp, and restart individual prepreg tows.\nIn the quest to automate the placement of composite materials at a high rate—to make the use of composites economical compared to conventional methods of fuselage fabrication—efforts have been focused at wrapping around a male mandrel, i.e. tool. Today's composite, fiber material placement processes and equipment have used male mandrels exclusively, wrapping tape layers on the outside surface of the tool—such as tool surface 108. One problem with this approach is that controlling the outside surface of the part—such as a fuselage section—is not possible without transferring the part to a female tool or clam shell type tooling.\nAs can be seen, there is a need for fabrication of composite parts using an automated lay-up machine that allows material placement directly to an outside mold surface, from inside a tool, allowing greater control and accuracy forming the exterior surface of the part. There is also a need for fabrication of composite parts using an automated lay up machine that eliminates expansion and transfer of the part to another tool, resulting in less defects and higher surface quality for the part. Moreover, there is a need for an automated lay-up machine for composite fabrication of large diameter fuselage sections."}
{"text": "Not Applicable\n1. Field of the Invention\nThis invention relates to drain strainers and stoppers. More particularly the present invention is directed to a combination flexible drain strainer and stopper which may be inserted into a drain opening and utilized to collect particulate matter while provide a plumbing fixture, such as a kitchen sink, with one or more drainage rates. The present invention is also directed to drain strainer capable of completely stopping drainage as desired.\n2. Description of the Related Art\nStrainers for use in sinks and other plumbing fixtures are common place in many homes and businesses. Such strainers are typically found in sinks, tubs, and other fixtures where they are utilized to prevent particulate matter such as food, hair and other potentially plumbing unfriendly items form passing out of the fixture and down the drain.\nMost strainers are somewhat permanently affixed above or within the drain opening of a plumbing fixture. While many strainers may be removed from the drain, removal may require that the strainer be pried out using a screw driver or other tool, or otherwise be removed with some difficulty. Where a sink includes a recessed or built in strainer, supplemental strainers have been developed to prevent particulate accumulation on such built in strainers, as such build up may accumulate and impede water flow.\nAn example of such a supplemental strainer is shown in U.S. Pat. No. 4,134,162 wherein a supplemental and disposable strainer is shown which has a conical straining surface which is secured over the drain opening and which projects upwardly above the bottom of the sink. While such a strainer is easily removed from the drain opening, the inherent shape, of the strainer makes it difficult to collect and dispose of particulate which may have accumulated thereabout. In addition, the strainer described in U.S. Pat. No. 4,134,162 does not disclose a means for adjusting the flow rate through the strainer or provide a plug which could be used to readily stop drainage through the strainer.\nWhile many sinks do not have a built in strainer within or over there drain openings. However, there are examples of upwardly projecting strainers for use with such drains. For example in U.S. Pat. No. 5,724,684 to Paar describes a raised strainer having a threaded region with a straining region mounted thereon. In Paar the threaded region may be threaded into a drain opening to allow the straining region to protrude upward into a sink. Paar also indicates that the use of a putty such as plumber\"\"s putty is needed to provide a water tight seal between the straining region and the sink surface. Again such a strainer does not provide for ready cessation of drain flow and does not provide for the ability to easily remove particulate which may accumulate thereon.\nSome other strainers have been developed which may be inserted into a drain basket so as to provide a removable strainer which may collect particulate therein. For example U.S. Pat. No. 3,982,289 to Robbins describes a disposable sink strainer which is shaped to fit complimentarily into a conventional sink strainer. While the Robbins sink strainer may easily be removed so as to dispose of particles which have accumulated therein, Robbins does not provide a strainer having more than one flow rate, or a means of stopping flow as is desired.\nThe present invention provides for a drain strainer that is removably securable to the bottom of a plumbing fixture, such as a sink, and which is constructed and arranged so as to encircle the drain opening therein. The present invention provides a fixture with a strainer capable of straining particulate matter out of a fluid flowing therethrough and which may be removed for cleaning and easy disposal of the particulate accumulated therein. As indicated, the present strainer may be provided with unique plug assembly which may be used to partially restrict the flow of fluid therethrough and which may be adjusted within the strainer to prevent flow altogether.\nThe present invention is made of a flexible material or materials such as rubber, polymer and/or polymer-like materials. The present strainer and plug are of relatively low cost and may easily be replaced if needed. In addition, the present strainer and plug may be installed and removed from a drain opening without the use of tools or sealants.\nIn light of the above, the present invention provides for a low cost sink strainer and an optional plug which may be inserted into a drain opening of any open plumbing fixture.\nIn at least one embodiment of the invention, the strainer may include an annular ring which provides the strainer with the ability to be self supporting within an open drain and form a seal around the drain to ensure that water flows through the strainer surface. Additionally, the annular ring provides for improved strainer capacity.\nIn at least one embodiment of the invention, the strainer may be utilized in conjunction with a plug to restrict or prevent water flow.\nIn at least one embodiment of the invention, the material is flexible and resistant to cutting, tearing, scratching and heat."}
{"text": "1. Field of the Invention\nThe present invention generally relates to protective devices and, more particularly, to a cushioning device for remote control television equipment.\n2. Prior Art\nVarious types of cushioning devices have been used with readily damagable equipment. For example, during the storing and shipment of clocks, watches, motor controls, electrical test equipment and the like, rigid foam polystyrene shells, strips and the like are used to encase and cushion such equipment against vibration and damage. Various configurations of such polystyrene protective containers are shown in U.S. Pat. Nos. 4,314,309, 4,602,715, 4,261,467, 4,201,293 and 3,346,221. Molded fiberglass has also been used for protective casings and the like. See, for example, U.S. Pat. No. 4,234,092.\nOnce delicate equipment arrives at its destination, however, it is customarily removed from its protective casing and is then used without any further protective containment. It has been found that remote control television equipment such as remote television channel selector panels, remote television on-off switches, remote videotape control panels and the like for television sets which are used without protective casings are especially subject to damage by being dropped or having solid objects fall in them as they rest on coffee tables, side and end tables and the like.\nIt would be desireable to be able to provide effective cushioning devices for remote control television equipment to encase the equipment while the equipment is in use, in order to reduce the possibility of their damage. Such cushioning devices should preferably be capable of protecting the equipment against vibrations, jars, drops, etc., and be adaptable to a wide variety of sizes and shapes of such equipment. Moreover, such devices should be inexpensive and durable."}
{"text": "1. Field of the Invention\nThe present invention relates to electrical switches, and particularly to an electromagnetic induction switch circuit.\n2. Description of Related Art\nAn ordinary laser pen includes a battery, a working circuit comprising a laser diode driver (LD driver) and a laser module, and a switch for selectively coupling the battery with the working circuit. Users can turn on and turn off the laser pen by operating a button, which connects to the switch.\nTherefore, users must manually press the button to turn on the laser pen when they want to use it, and have to press the button again to turn off the laser pen when not using it, to save power. It is inconvenient for users to operate the switch when they use the laser pen frequently. Oftentimes, users forget to turn off the laser pen when not using it, which results in waste of power.\nTherefore, an electromagnetic induction switch circuit which can automatically turn a laser pen on or off is needed."}
{"text": "Traditionally, chlorofluorocarbons (CFCs) like trichlorofluoromethane and dichlorodifluoromethane have been used as refrigerants, blowing agents and diluents for gaseous sterilization. In recent years, there has been universal concern that completely halogenated chlorofluorocarbons might be detrimental to the Earth's ozone layer. Therefore, stratospherically safer alternatives to these materials are desirable.\nConsequently, there is a worldwide effort to use fluorine-substituted hydrocarbons which contain fewer or no chlorine substituents. The production of HFCs, i.e., compounds containing only carbon, hydrogen and fluorine, has been the subject of interest to provide environmentally desirable products for use as solvents, blowing agents, refrigerants, cleaning agents, aerosol propellants, heat transfer media, dielectrics, fire extinguishing compositions and power cycle working fluids. It is known in the art to produce fluorocarbons such as HFCs by reacting hydrogen fluoride with various hydrochlorocarbon compounds. Such HFCs are not only considered to be much more environmentally advantageous than hydrochlorofluorocarbons (HCFCs) or chlorofluorocarbons (CFCs) because they are not non-ozone depleting, but also they are non-flammable, and non-toxic as compared to the chlorine containing compounds. While HFCs are considered to be much more environmentally advantageous than HCFCs or CFCs because they are non-ozone depleting, recent data indicates that they may also contribute to greenhouse global warming. Accordingly, alternatives to HFCs, HCFCs, and CFCs are also being explored.\nHydrofluoroolefins (“HFOs”) have been proposed as possible replacements. It is generally known that HFOs are best used as a single component fluid or azeotropic mixture, neither of which fractionate on boiling and evaporation. The identification of such compositions is difficult due, at least in part, to the relative unpredictability of azeotrope formation. Therefore, industry is continually seeking new HFO-based mixtures that are acceptable and environmentally safer substitutes for CFCs, HCFCs, and HFCs.\n3,3,3-trifluoropropyne, also known as “TFPY”, is a suitable starting material for various hydrofluoroolefins (HFOs) including but not limited to HFO-1234yf, HFO-1234ze, HCFO-1233zd(Z), and HFO-1225ye, and is also produced as a side product in the production of various hydrofluoroolefins including but not limited to HFO-1234ze, HFO-1234yf, HFCO-1233zd, all of which are well known in the art, and several are described in U.S. Patent Application Publication No. 2009/0234165, assigned to the assignee of the present invention, the disclosure of which is incorporated herein by reference.\nMethods to produce 3,3,3-trifluoropropyne are also known in the art, for example, from U.S. Pat. Nos. 7,964,759 and 8,791,309.\nNew compositions of, methods of separating and purifying, and uses of, 3,3,3-trifluoropropyne are desired."}
{"text": "This invention relates to vehicular disc brakes of the sliding caliper type and, more particularly, to a novel torque plate for use with a sliding caliper disc brake.\nIn a typical sliding caliper disc brake, the caliper embraces the rotor and is slidably supported at its opposite ends by a torque plate fixedly secured to a non-rotating portion of the axle assembly of the associated vehicle. As the brake is applied by the actuator assembly, the brake pad driven directly by the actuator assembly, typically the inboard pad, is pressed against the inboard face of the rotor, whereafter, with continued driving input from the actuator assembly, the caliper slides inboard on the torque plate to bring the outboard pad, carried by the caliper, into frictional engagement with the outboard face of the rotor. The torque plate is a critical element in this combination. In addition to allowing smooth sliding movement of the caliper on the torque plate, the torque plate, particularly in heavy-duty applications, must efficiently absorb large magnitudes of braking torque and transmit this torque efficiently to the fixed vehicle assembly, typically the axle housing. The torque plate must also be configured and installed in such a manner as to not significantly interfere with the flow of cooling air over the rotor. Prior art torque plates have either sacrificed brake cooling effectiveness to maximize torque absorbing capacity, or have limited torque absorbing capacity to maximize brake cooling efficiency."}
{"text": "Early technology for vascular access required surgical opening of the skin, viewing of the vessel to be catheterized, and placement, under sterile procedure of a stainless steel needle. After the insertion process, the skin was approximated and the access site covered to prevent infection. The process required a competent physician, and was indicated only in rare instances when the patient required intravascular sustenance or required replacement fluids in association with major surgery.\nWith many of the described catheterization technologies, blood contamination was a risk when the intra-vascular catheter was inserted. Blood flashback often escaped the catheter hub when the guide needle was removed and before the intra-vascular solution tubing could be connected, thus exposing the caregiver and patient to blood leakage. Several notable valve designs have been patented to reduce this blood leakage. For example, in U.S. Pat. No. 4,387,879, Tauschinski describes a self-sealing elastomeric disc that can be incorporated into a connector body to interface with a parenteral supply solution and an intra-vascular catheter. Similarly, Motisi et al, describe a one way valve in the body of a catheter apparatus in U.S. Pat. No. 5,843,046. However, the necessity of a plunger and introduction of needle or tubing through the plunger decreases the inside diameter of the catheter and reduces the fluid flow rates for this design. An “O” ring in the valve can prevent leakage from around the plunger, but this decreases the inside diameter of the large bore catheter. The valve is held in place by a “cap” which puts its placement deeper into the throughbore, out of reach of conventional intra-vascular tubing or conventional syringes. The “cap” also prevents connection to conventional tubing or conventional syringes. This valve is bulky and decreases the size of the intra-vascular catheter. It also cannot be opened by insertion of conventional devices used by those skilled in the art.\nI-V Access Technology, Inc. (IVAT) has intellectual property for a peripheral IV catheter that incorporates an automatic needle retraction mechanism (U.S. Pat. Nos. 8,105,288 and 8,591,469) to retract or withdraw the guide needle when the vein is accessed, a dilator to expand the lumen created by the guide needle (also in US patent divisional application Quine docket 138-000112US), and a pressure activated fluid flow control valve (US patent application US-2013-0204226-A1) to prevent blood leakage when the needle or needle/dilator assembly is fully withdrawn from the catheter assembly. The valve is opened and IV therapy administered through the catheter and into the vasculature when a male luer fitting is inserted into the proximal opening of the catheter hub and contacts the proximal side of the valve.\nAlthough the specific design of the male luer fitting is standardized by ISO, in practice the manufacturers of IV fittings, luer connectors, and syringes produce their products with the male luer geometry sometimes deviating significantly from the ISO standard. As a result, the degree to which the male luer fitting contacts the proximal side of IVAT' s valve can vary functionally from an inability to open the valve (no contact at all) to damaging the valve from a shearing of the valve against the catheter hub distal to the valve.\nOther manufacturers who produce IV catheters with integrated valves have incorporated a separate slidable adapter piece within the catheter hub that, as the male luer is inserted, it contacts the adapter piece which is axially displaced and penetrates the slits in the valve, thereby opening the fluid flow pathway. The adapter piece that penetrates the slits is generally conical in shape, which allows the adapter to readily open the valve without excessive force, and, when the male luer is removed (e.g. when medication injection through a syringe is completed) the valve exerts an elastic force on the adapter to axially displace the adapter and close the valve to fluid flow. The adapter piece, while considered a necessity to accommodate various male luer geometries and dimensional variations, adds cost and complexity to the manufacturing of the IV catheter. The extra component requires extra material, tooling, and labor costs to produce and potentially additional costs for assembly.\nIn light of the problems remaining in the art, it would be beneficial to have catheter devices well adapted to receive both standard and non-standard male luer fittings. It would be desirable to have sanitary catheter fittings that seal with non-standard male luers while allowing automatic opening of valves facilitating withdrawal or injection of fluids. Simplified methods and component configurations are needed to allow ready assembly of catheter devices containing a variety of elements finely crafted from a variety of materials. The present invention provides these and other features that will be apparent upon review of the following."}
{"text": "1. Field\nOne or more exemplary embodiments relate to a connector and a method of manufacturing the same.\n2. Description of the Related Art\nRecent advancements and refinements in the semiconductor industry have shrunk sizes of semiconductor devices. Due to their shrunken sizes, the semiconductor devices have become more vulnerable to an influence of electrostatic discharge (ESD) or a surge. Abnormalities caused by ESD can occur not only in semiconductor manufacturing processes but also in electronic parts and semiconductor devices manufactured by such semiconductor manufacturing processes.\nA method of embedding a diode has been applied to prevent ESD in an integrated circuit (IC). However, an embedded diode may decrease a performance of the IC. To prevent a semiconductor device from ESD, a circuit including a transient voltage suppressor (TVS) diode has been used."}
{"text": "The present application is based on and claims priority from Japanese Patent Application 2002-9931, filed Jan. 18, 2002, the contents of which are incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to an ac generator that has a rotor having a plurality of claw-shaped pole pieces and a field coil held inside the pole pieces.\n2. Description of the Related Art\nIt has been found that the magnetic noise of an ac generator is caused by the third harmonic wave of magnetic flux that is distorted during operation. In order to reshape the waveform of the magnetic flux into a sine wave, JP-B2-61-11066 proposes an ac generator in which a rotor has scalene trapezoidal claw-shaped pole pieces. That is, the surface of the pole pieces has a trapezoidal shape that is asymmetrical with respect to the center line of the pole piece that is perpendicular to the rotation direction of the rotor.\nJP-B2-3223536 proposes a rotor having claw-shaped pole pieces each of which has a chamfered front surface that is formed at the front of the pole pieces in the rotation direction thereof and a chamfered rear surface that is formed at the rear of the pole pieces in the rotation direction and is larger than the chamfered front surface.\nHowever, it is difficult to form the claw-shaped pole pieces disclosed in JP-B2-61-11066 by forging because they are asymmetrical with respect to the center line, which makes it difficult to evenly apply forging pressure on the entire surface of the pole pieces. As a result, the rotor can not have the pole pieces having an accurate uniform shape.\nOn the other hand, the ac generator disclosed in JP-B2-3223536 has a problem of a magnetic noise caused by sharp change of magnetic flux between the pole pieces. Because the chamfered rear surface of one pole piece and the chamfered front surface of the next pole piece extend from opposite end of the rotor to face each other, the magnetic flux changes sharply when the rotor rotates, thereby causing magnetic noise. In addition, the output power may be reduced because the chamfered rear surface becomes too large for the pole pieces to conduct effective magnetic flux.\nTherefore, an object of the invention is to provide a powerful and noiseless ac generator.\nAn ac generator according to a feature of the invention includes a rotor having trapezoidal claw-shaped pole pieces that alternately extend in opposite directions to interleave each other, in which each of the pole pieces has a chamfered surface at the tip thereof that faces the inside surface of the stator. The chamfered surface has narrower area at the front half of the pole piece in the rotation direction than the rear half. Because such an asymmetrical chamfered surface is formed at limited area after the pole pieces are forged, the pole pieces having accurate size and shape can be provided. When the rotor with such pole pieces and a field coil is rotated in a predetermined direction to generate power, the magnetic flux sharply increases at the front half of the chamfered surfaces of the pole pieces. On the other hand, the magnetic flux gradually decreases at the rear half of the chamfered surfaces. Therefore, a magnetic wave that is similar to the wave provided by the scalene trapezoidal pole pieces of a prior art ac generator is formed. Such a magnetic wave is opposite in phase to the third harmonic wave, so that the magnetic noise caused by the third harmonic wave is canceled. In addition, because the area of the chamfered surface is comparatively smaller than the prior art, reduction in the output power is negligibly small.\nIn an ac generator according to another feature of the invention, a rear half surface of each of the pole pieces in the rotation direction has a concavity or a through hole at the tip portion thereof facing the inside surface of the stator core. The concavity or the through hole provides an asymmetrical chamfered surface at a limited area of the pole piece. Because the concavity or the through hole can be easily formed when or after the pole pieces are forged, the pole pieces having accurate size and shape can be provided. This feature brings about substantially the same effects as the above-described ac generator.\nIn the ac generator described above, each of the pole pieces may have a front chamfered surface and a rear chamfered surface that are approximately symmetrical at the opposite corners thereof. Therefore, there is little change in magnetic flux between the pole pieces. As a result, the magnetic noise caused by change in the magnetic flux is very low."}
{"text": "1. Field of the Invention\nThis invention relates to an improved control system for maintaining a substantially constant mass flow of wire through the successive stations of a continuous cumulative wire drawing machine.\n2. Description of the Prior Art\nContinuous cumulative wire drawing machines can be classified as being either of the \"single block\" or \"double block\" types. In single block machines, as the term implies, at each drawing station the wire is drawn through a die by a single rotatably driven block. A segment of the wire is allowed to accumulate as multiple windings on the block before passing through a guide in the form of a pay-off eye on an overlying carrier ring which is freely rotatable about the block axis. From here, the wire continues over appropriately arranged guide sheaves either to another die at the next drawing station or to the final take-up block.\nIn double block machines, upper and lower blocks are arranged at each drawing station on a vertically disposed drive shaft. The lower block is keyed to the drive shaft for rotation therewith, and serves to draw the wire through the die. The upper block is freely rotatable on the drive shaft, and an independently rotatable carrier ring with a guide consisting of a transfer sheave is interposed between the two blocks. After passing through the die, wire accumulates on the lower block before passing over the transfer sheave onto the upper block where a further accumulation takes place. The wire then continues from the upper block over other appropriately positioned guide sheaves to the next station.\nUnder an ideal constant mass flow condition, i.e., when the mass flow rate at which wire is being wound onto the blocks equals the mass flow rate at which wire is being payed off the blocks, the carrier rings and their respective pay-off eyes or transfer sheaves will remain stationary. More often as not, however because of uneven die wear and/or other variable operating conditions, these rates will differ at one or more of the drawing stations. The carrier rings will compensate for differences between take-up and pay-off rates by rotating in either clockwise or counterclockwise directions, depending on whether the accumulations of wire on the blocks are increasing or decreasing.\nThere is, however, a limit to the extent to which such increases or decreases in wire accumulations can be tolerated. Thus, operating personnel must constantly monitor and adjust the wire accumulation on each block. Where the blocks are all connected via clutches to a common drive, as is usually the case, this entails frequent disengagement and reengagement of the clutches of selected blocks. This is a burdensome task, and one that requires considerable experience and a high degree of skill.\nMoreover, when a clutch is disengaged to momentarily stop a selected block, the cooling time for the wire segment accumulated on that block will be extended as compared with the cooling times of the wire segments passing around those other blocks which continue to rotate. Any such localized extended cooling may produce a localized unacceptable variation in the metallurgical properties of the wire.\nIn the past, attempts have been made at alleviating the control responsibilities of operating personnel by providing systems designed to automatically monitor and adjust mass flow conditions. Typically, such systems employ vertically spaced pairs of photoelectric cells and light sources arranged to define the upper and lower limits of wire accumulation on the blocks. The signals generated by the photoelectric cells are used to automatically engage and disengage the clutches connecting the blocks to the common drive. While such systems can relieve operating personnel of some control responsibility, they still produce unacceptable localized variations in metallurgical properties caused by intermittent operation of the blocks. Moreover, such systems often malfunction as a result of the photoelectric cells becoming coated with the dust which usually pervades the atmosphere of a wire mill.\nAn object of the present invention is to provide an improved control system for maintaining a substantially constant mass flow rate of wire through the successive stations of a continuous cumulative wire drawing machine.\nAnother object of the present invention is to achieve the aforesaid substantially constant mass flow rate without imparting undesirable localized variations to the metallurgical properties of the wire as a result of uneven cooling.\nStill another object of the present invention is the provision of a control system which can operate reliably in dust laden environments of the type often found in wire drawing mills."}
{"text": "1. Field of the Invention\nThis invention relates to solid-electrolyte fuel cell electrode material, and to an electrode formed of the material.\n2. Description of the Prior Art\nConventionally, solid-electrolyte fuel cell electrodes are usually formed of LaMnO.sub.3, LaCoO.sub.3, LaSrMnO.sub.3 or other such lanthanum-based material having good heat-resistance and a coefficient of thermal expansion that is close to that of the solid electrolyte material. To manufacture electrodes using such materials, the surface of a solid electrolyte such as stabilized zirconia is flame-coated with a powder of the electrode material, or coated with a slurry of the powder which is then sintered, producing an air electrode in the form of a porous membrane having good air permeability.\nWhile such air electrodes are required to have good air permeability, they are also required to have good electrical conductivity. The conductivity of an air electrode is highly dependent on its air permeability, but higher permeability means decreased density, with a corresponding decrease in conductivity; that is, the electrical resistance is increased.\nAn object of the present invention is to provide a solid-electrolyte fuel cell electrode which provides a major improvement in conductivity without a reduction in air permeability."}
{"text": "The invention relates generally to a copier.\nMore particularly, the invention relates to a photographic copier, especially a copier of the type in which the originals and the copying material are supplied from rolls.\nA known photographic copier of this type is designed to copy photographic originals, e.g. negatives, having different widths onto bands of copying material which likewise have different widths. The copier is provided with guide means for the originals as well as the copying material. Each guide means may consist of a set of interchangeable guides which are dimensioned for different widths. Alternatively, each guide means may consist of a single guide which is movably mounted in the copier so that it may assume different positions for different widths. The copier further has an adjustable objective, as well as means for regulating the lengths of the images of the originals.\nSuch copiers, which find widespread application, are marketed in various forms. As a rule, adjustment of a copier of this type assumes a knowledge of which copying material is to be used for originals having given dimensions. Consequently, when the copier is to be adjusted for a particular original, the required width of the copying material or paper is known, as are the dimensions of the original which constitute an important parameter for photoelectric scanning of the original. Also known is the length of copying material to be exposed, and the cross section of the reflector shaft or tube which will provide optimum illumination of the original. Furthermore, the magnification desired from the objective is known. In this regard, either of two types of adjustable objectives may be employed. One type of objective has a variable focal length. The other type has a fixed focal length and is mounted for movement towards and away from the original and, hence, towards and away from the copying material.\nThe number of adjustments is generally large since an adjustment is usually required for each of the parameters outlined above. In a simple copier, all of the adjustments are performed manually which carries with it the risk of an error. This risk is compounded by the fact that the error may occur in the darkened portion of the copier which cannot be observed during operation.\nVarious copiers are known in which some of the adjustments are perform automatically. For instance, the German Pat. No. 17 97 492, which corresponds to the U.S. Pat. No. 3,488,779 discloses a copier which is designed in such a manner that the intensity of illumination, the compensation for the size of the original, the mask for the copying material, and the length of copying material to be exposed are all adjusted automatically via suitable adjusting elements in response to installation of a mask for the original. However, conventional copiers having motorized adjusting elements are extremely expensive. In this regard, it is frequently cheaper to provide for interchange of relatively simple parts, instead of motorized adjustment of such parts, in order to adapt a copier to originals of different sizes. The construction of the known motorized copiers does not, however, permit the interchange of parts."}
{"text": "Non-invasive (in vivo) methods for measuring various blood-related parameters have become very popular due to the fact that these measurements, in distinction to invasive ones, do not involve the physical withdrawal of a blood sample from the patient's body. Optical monitoring techniques of the kind specified utilize the detection of light transmitted or reflected from the location on the patient's body under measurement, and are based on spectrophotometric measurements enabling the indication of the presence of various blood constituents based on known spectral behaviors of these constituents. These methods being applied in real medicine rather than in analytical chemistry create the basis for non-invasive blood tests, which present, no doubt, one of today's most exciting challenges. To make blood tests low-cost, safe and painless means to make them non-invasive.\nThe two main challenges, that any non-invasive optical method has to deal with, are as follows: (1) the low signal-to-noise ratio, and, (2) the large variability of individual parameters influencing the signal of concrete patients.\nMost of these techniques utilize a measurement optical device or probe, designed in a manner to be attached to the patient's finger, which includes an optical assembly for irradiating the finger with light and detecting its light response. The conventional devices of the kind specified, such as a pulse oximeter, which is the generally accepted standard of everyday clinical practice, provide for measuring enhanced optical pulsatile signals caused by the changes in the volume of a blood flowing through a fleshy medium (e.g., finger).\nIt is known that for blood parameters other than oxygen saturation, e.g., glucose concentration, significant difficulties have been encountered, because their absorption spectral behavior in red and near infrared regions is not as remarkable as for the oxygenized hemoglobin. Hence, the main limitations on the way of expanding the non-invasive techniques to the measurements different from pulse oximetry are associated with the limited selectivity of the absorption based method.\nA different technique is disclosed in U.S. Pat. No. 6,400,972, WO 01/45553 and WO 01/96872, all assigned to the assignee of the present application. This is an occlusion-release based technique, according to which an over-systolic pressure is applied to the blood perfused fleshy medium with a normal blood flow so as to create a state of temporary blood flow cessation at the measurement location. The measurement with different wavelengths of incident radiation and/or different polarization states of detected light are carried out at timely separated sessions taken during a time period including a cessation time when the state of the blood flow cessation is maintained. This technique utilizes the condition of the “artificial blood kinetics” rather than the natural blood kinetics taking place when the state of blood cessation is not achieved. As a result of the cessation of the blood flow, a condition of the artificial kinetics is achieved with the optical characteristics of the blood associated with the light response being different from those at the natural blood kinetics. Indeed, it is known that the scattering properties of blood depend on the size and shape of scatterers (aggregates). Thus, time changes of the light response at the condition of artificial kinetics depend on the changes in the shape and average size of the scattering centers in the medium, i.e., red blood cells (RBC) aggregation (Rouleaux effect). It was found that owing to the effect of the artificial kinetics, the optical characteristics of blood changes dramatically, such that they differ from those of the fleshy medium with a normal blood flow by about 25 to 60%, and sometimes even more. Hence, the accuracy (i.e., signal-to-noise ratio) of the technique based on the artificial kinetics as well as selectivity of the optical measurements can be substantially better when compared with those based on measurements of the blood parameters at natural kinetics."}
{"text": "It is often desirable to plant large numbers of genetically identical plants that have been selected to have advantageous properties, but in many cases it is not feasible to produce such plants using standard breeding techniques. In vitro culture of somatic or zygotic plant embryos can be used to produce large numbers of genetically identical embryos that have the capacity to develop into normal plants. However, the resulting embryos lack the protective and nutritive structures found in natural botanic seeds that shelter the plant embryo inside the seed from the harsh soil environment and nurture the embryo during the critical stages of sowing and germination. Attempts have been made to provide such protective and nutritive structures by using manufactured seeds, but so far germination from manufactured seeds is less successful than from natural seeds. There remain large differences between manufactured seeds and corresponding natural seeds. Whereas the embryo relies on the megagametophyte for nutrients useful for germination, the embryo in a manufactured seed relies on the nutritive medium that is provided in the manufactured seed.\nTherefore, there is a need for generating improved nutritive medium that is useful for improving rates of conversion for manufactured seeds containing somatic embryos to provide a large number of normal germinants. The present invention addresses this and other needs."}
{"text": "1. Field of the Invention\nThe present invention relates to a tillage system used by farmers to till a field in spring prior to planting and in fall after harvesting farm crops; and, more particularly to a tillage system used to produce deep tilled rows and shallow tilled regions therebetween, leaving behind a smooth clean-tilled field that conserves water and soil fertility without excessive soil erosion.\n2. Description of the Prior Art\nMany patents address issues related to tilling a farm field. Several of the tillage systems are designed to deep-till a field. Such deep-tilling causes top soil loss due to soil erosion. In addition, deep-tilling degrades field fertility more and more each year. Run-off of soil and fertilizer tends to pollute rivers and waterways. Some of the patents restrict deep tilling to spaced rows. These tillage systems practice minimal tillage principles and leave behind compacted soil between deep tilled rows with stocks of weed protruding from the soil resulting in unsightly field after the tillage operation. Besides, the compacted soil between the deep tilled rows encourages water run offs and soil erosion at the deep tilled rows.\nU.S. Pat. No. 3,941,193 to Shoemaker discloses a ground tilling apparatus. The tilling apparatus characterized by a regulating drive connection between the towing tractor drive axle and the tiller shaft, the tiller tines having a modified arcuate configuration which provide a pocketed, rather than a furrowed or channeled, tilling depth profile. The tiller shaft drives the tiller in the same direction with the tiller circumferential speed substantially equal to the ground speed of the tractor, whereby said tips generate a curtate cycloidal curve as the tines rotate and the tips pierce the ground by curtate cycloidal movement with the shank portions following in the holes formed by the tips. This tilling apparatus has two sets of tilling tines oriented and turned in opposite directions. The first tine rotational speed is matched to the ground speed so as to remove large clumps of soil and the second set of tines bust the clumps. The device does not till deep along planting rows.\nU.S. Pat. No. 4,102,406 to Orthman discloses a ground conditioning device and method of conditioning soil. This device creates water-holding and retaining cavities that are formed in furrows with the cavities being staggered between adjacent rows. An arm spring biased downwardly carries a hub having spikes in planes laterally spaced apart with the spikes around the hub being staggered between the planes. Each spike is concave longitudinally and in transverse cross section and is secured to opposite sides of a disc mounting plate by a pair of clamping elements one of which is integrally connected to the mounting plate. The hiller row unit is used after the row crop is established in a bed prepared by the splitter or shaper equipment that is disclosed in U.S. Pat. Nos. 3,718,103 and 3,870,108. The hiller provides for ridging with an absolute minimum of root pruning with only the soil that is used to build the ridge being disturbed. The device does not deep-till a plurality of parallel rows where seed planting is planned.\nU.S. Pat. No. 4,834,189 to Peterson et al. discloses a row crop cultivator. The row crop cultivator is adapted for use in minimum tillage applications. A plurality of ganged cultivating units is mounted on a tool bar adapted to be towed behind a tractor vehicle. Each cultivating unit comprises a frame, which is supported on a pair of spaced-apart gauge wheels. Also mounted on the frame for rotation within the gap between the gauge wheels is a disk-shaped coulter blade, which is configured to cut through crop residue and weed debris as the material being cut is held against the ground surface by the gauge wheels. This arrangement tends to prevent improper cutting known as “hair-pinning”. Trailing directly behind the coulter blade is a middleworker comprising a narrow width shank depending from the frame. The shank has a weld point on its lower leading edge and wing-like share blades flaring rearwardly and outwardly at a predetermined angle from the lower end of the shank. The middleworker design reduces the tendency for slabs of earth to be thrown onto and thereby damaging growing row crops. The row crop cultivator shields the row of plants that are growing in rows with a pair of coulters and the portion between the rows is weeded or disturbed by a horizontal tiller. The row crop cultivator does not deep-till the field at row planting locations.\nU.S. Pat. No. 5,623,997 to Rawson et al. discloses a soil zone-builder coulter closer/tiller. The soil manipulation device is provided for closing a groove in soil formed by a soil tiller shank. The device includes a frame connected to a rear portion of the tiller shank and a pair of rotatable coulter blades mounted on the frame for engaging soil behind the tiller shank. The coulter blades are positioned on opposite sides of the groove so as to close the groove in the soil behind the tiller shank. The coulter blades are selectively movable and positionable in directions toward and away from each other so as to vary the distance between the coulter blades. The coulter blades also are selectively movable and positionable in a substantially vertical direction with respect to the frame, so as to permit variation in depth of coulter blade penetration into the soil. The coulter blades have side surfaces, which are selectively movable and positionable from orientations which are substantially in parallel with the groove to orientations which are out of parallel with the groove, so as to variably manipulate soil on opposite sides of the groove. The tiller meant for inserting seeds is followed by two coulters that close the soil groove formed by the tiller. The coulters may be independently positioned to adjust the groove closing process. The device of the '997 patent does not till a field; but rather plants seeds in a previously tilled field and closes the grove formed after setting plant seeds in the groove.\nU.S. Pat. No. 7,575,066 to Bauer discloses a zone tillage tool and method. This patented device is marketed as “Soilwarrior”, and detailed by the web page at http://www.soilwarrior.com/soil.php. This zone tillage tool comprises a tool frame connected by a parallel lift linkage to a tool bar. A fall tillage coulter or dual spring tillage coulters can be carried on the tool frame in advance of a pair of containment coulters. The fall tillage coulter tills soil in a zone that is approximately 7 to 9 inches deep while the spring tillage coulters till 2 to 3 inches deep. The tillage tool has a tool bar that mounts a plurality of tools with its own tool frames carrying independently controllable tillage coulter followed by a pair of slanted containment coulters in the direction of movement of the tillage tool. The tillage tool with its plurality of tillage coulter and containment coulters only creates deep tilled grooves in the field at the places where planting will be done. Spaces between these deep grooves are not tilled, and the soil remains undisturbed with all the weeds still in place, creating an unsightly appearance, Moreover, when rain water or irrigation water is applied to the field, all the water runs into the deep grooves, which have only limited volume and ability to absorb water. This causes selective erosion of soil from the deep grooves. The rainwater or irrigation water runs freely on the compacted soil between the deep tilled rows; run-off water rich in fertilizers and top soil pollutes lakes and rivers.\nA number of prior art patents disclose coulters of different geometry and harrowing tools. For example, U.S. Pat. No. 369,163 to Clark discloses a disk harrow: U.S. Pat. No. 1,954,783 to Bohmker discloses a disk blade; U.S. Pat. No. 4,538,688 to Szucs et al, discloses fluted coulter blade and U.S. Pat. No. 5,649,602 to Bruce discloses wavy coulter.\nThere remains a need in the art for a tillage system that deep-tills at the planting rows so that planted seeds grow easily and vigorously. The entire field should not be deep-tilled, since deep-tilling everywhere results in water run-offs that produce soil erosion and fertilizer loss with consequent pollution of rivers and water ways. There is need in the art for a tilling system that clears plant debris, including stocks, sufficiently to leave behind a field with a clean-tilled appearance and minimal soil disturbance between the deep-tilled rows and which provides sufficient water absorption and storage to prevent water run-offs and consequent soil erosion. There is a further need in the art for a tillage system that can be readily added-on to existing devices and tools."}
{"text": "Commercial and recreational fishing is often conducted in shallow water. Both fresh water and salt water shallows are often populated with a variety of fish. Fishermen who fish these waters precisely locate and anchor their boat in areas where the amount of fish caught will be maximized. Often times fishermen will locate their boat where fish are visually detectable within the water through a technique called sight fishing. When using this technique, the fishermen must make every attempt to minimize noise so as not to scare the fish.\nConventional anchors are typically used to anchor a boat when shallow water fishing. A conventional anchor may take on various forms but generally has the form of a mass located at the end of a rope or chain that is in turn attached to the boat. To anchor the boat, a fisherman simply drops or throws the mass into the body of water letting it sink to a bottom thereof.\nUnfortunately, several problems arise when using a conventional anchor during shallow water fishing. First, because the anchor is ordinarily tethered to the boat using a rope or chain, the boat will drift when anchored due to currents within the water. This drifting effect can place the boat in an unintended position other than a position most advantageous for shallow water fishing. Second, a loud noise and splash is produced when the anchor is thrown into the water that in turn can scare away the fish in proximity to the boat. Third, the mass often times drags across the bottom surface of the body of water and stirs up particulate matter such that the fisherman's view of fish within the water is obscured. Additionally, the mass can damage the vegetation growing at the bottom of the body of water as it drags across it.\nIn view of the above, it is desirable to have an anchor that anchors a watercraft within the water such that the watercraft does not drift due to current. It is further desirable that such an anchor function without producing an excessive amount of noise or obscuring the clarity of the water.\nEmbodiments of the invention provide such an anchor. These and other advantages of the invention, as well as additional inventive features, will be apparent from the description of the invention provided herein."}
{"text": "Manganese (IV) compounds are well known and are used in a variety of oxidation reactions. For example, manganese dioxide (MnO.sub.2) has been used in the manufacture of chlorine gas from hydrogen chloride and the oxidation of aniline to hydroquinone. See \"Chemistry of the Elements\", N. N. Woodward and A. Earnnshaw, Pergammon Press, Oxford, pp 1219-20 (1984). A molecular manganese-oxo cluster is involved in the oxidation of water to oxygen in the photosynthesis process used by plants. See Yachandra et.al., Science, 260, 675-679 (1993). Because manganese has stable oxidation states of +4, +3 and +2, manganese oxides can be used in batteries.\nManganese oxides can have layered structures or three-dimensional microporous structures. S. Bark et al., Electrochimica Acta, 36, 1595-1603 (1991), P. LeGoff et al., Mat. Res. Bull., 31, 63-75 (1996), P. Strobel et al., Mat. Res. Bull., 28, 93-100 (1993), Y. Shen et al., Science, 260, 511-515 (1993). Finally, the ion-exchange properties of manganese oxide compositions have been reported by Q. Feng et al. in Chem. Mater., 7, 148-153 and 1722-1727 (1995).\nIn addition to manganese oxide compounds, there are reports of manganese phosphate molecular complexes containing Mn(IV). For example, Weighardt et al. have reported the synthesis of a trimeric manganese (IV) phosphate complex. Weighardt et al., Chem. Comm., 1145 (1988). A Mn(IV) dimer having the formula [(bpy)(H.sub.2 PO.sub.4)Mn.sup.+4 (.mu.-O).sub.2 (.mu.-HPO.sub.4)Mn.sup.+4 (bpy)]*H.sub.2 O where bpy=bipyridine has been reported by J. Sarneski et al., Inorg. Chem., 30, 2833-35 (1991). M. Jaky has disclosed dilute manganese (IV) phosphate solutions which were used for the oxidation of various organic substrates. Polyhedron, 12, 1271 (1993). Finally, there are a number of examples of Mn(III)-containing phosphates that have been prepared by hydrothermal synthesis, e.g., KMn.sub.2 O(PO.sub.4)(HPO.sub.4), Lightfoot et. al., J. Solid State Chem., 73, 325-329, (1988), and NH.sub.4 Mn.sub.2 O(PO.sub.4)(HPO.sub.4), Lightfoot et. al., J. Solid State Chem., 78, 17-22, (1989).\nIn contrast to these references, applicants have synthesized crystalline manganese phosphate compounds which contain Mn(IV) and which have an extended network. By extended network is meant that the defining Mn--P--O structural unit of the material repeats itself into at least two adjacent unit cells without termination of bonding, i.e., the material is not molecular. See \"Structural Inorganic Chemistry, Fifth Edition,\" A. F. Wells, Clarendon Press, Oxford, pp. 11-15, (1984). The network can be one-dimensional (a linear chain), two-dimensional (layered) or three-dimensional. The three dimensional network may or may not be a microporous network. There is no mention of a manganese(IV)-containing phosphate material with an extended network. By Mn(IV)-containing phosphate, it is meant that the average oxidation state of Mn is greater than 3.0 but less than or equal to 4.0, indicating the presence of some Mn(IV). These compositions are prepared by employing large excesses of phosphate to avoid the precipitation of insoluble Mn(IV) oxides, and through careful control of the pH and the use of mild reaction conditions. Further, applicant has also synthesized metallo manganese phosphates where a portion of the manganese is replaced by a metal such as iron (III), aluminum, gallium, etc. This new family of compounds are useful in several hydrocarbon conversion processes such as oxidative dehydrogenation of alkanes and oxidative coupling of methane to higher hydrocarbons."}
{"text": "1. Technical Field\nThe present invention relates to a Surface Acoustic Wave (hereinafter called SAW) filter used in a mobile phone or the like, an antenna duplexer using the same, and a mobile communication terminal using the same.\n2. Prior Art\nGenerally, the SAW filter used in the radio frequency spectrum such as used by a mobile phone typically includes a ladder-type circuit to realize a band-pass filter of small size and high selectivity. Hitherto, these types of ladder-type SAW filter was constituted as shown in FIG. 17 as disclosed in Japanese Laid-open Patent No. 9-102728.\nThat is, serial-branch SAW resonators 12-1 to 12-4 include reflectors 15 sandwiching an interdigital transducer 14 are arranged at the left side (the side of L0) in the drawing on a piezoelectric substrate 13, and are connected by serial-branch wiring patterns 16-1 to 16-3. On the other hand, parallel-branch SAW resonators 12-5 to 12-7 including reflectors 15 sandwiching an interdigital transducer 14 are arranged at the right side (the side of L1) and are connected to the serial-branch wiring patterns 16-1 to 16-3 by parallel-branch wiring patterns 16-4 to 16-6.\nIn this constitution, a seven-element ladder-type SAW filter 11 is formed. This SAW filter 11 has a circuit configuration in which the input electrode 16-8 side starts from the serial-branch SAW resonator 12-1. Such circuit configuration is called the serial-branch input type. By contrast, the circuit configuration in which the input side starts from the parallel-branch SAW resonator is called the parallel-branch input type.\nThe conventional ladder-type SAW filter 11 tends to be less in increase of reflection coefficient at the input electrode 16-8 even outside of the pass band. FIG. 18 is a polar coordinate chart of reflection coefficient as seen from the input electrode 16-8 of a receiving filter e.g. a receiving filter for Advance Mobile Phone Service (hereinafter called AMPS) which is a mobile phone system in the United States. The receiving band, that is, the pass-band is 869 MHz to 894 MHz (marker M1 to M2), and the transmitting band to be rejected at the receiving input terminal is 824 MHz to 849 MHz (marker M3 to M4).\nAs known from FIG. 18, at the upper end M4 of the transmitting band, the reflection coefficient of only about 0.8 is obtained. The cause is the inductive components of the parallel-branch wiring patterns 16-4 to 16-6 as explained in detail later in the embodiment. In the case of serial-branch input type, in particular, since the SAW resonators 12-1 to 12-7 are arranged as shown in FIG. 17 for the convenience of constitution, the parallel-branch wiring pattern is extended in length, and this tendency appears obviously.\nThe use of this ladder-type SAW filter 11 in an antenna duplexer is explained below. The antenna duplexer is a branching filter comprising a transmitting terminal, an antenna terminal, and a receiving terminal, and a transmitting filter is placed between the transmitting terminal and the antenna terminal, and a receiving filter is placed between the antenna terminal and the receiving terminal.\nMoreover, phase shifters are provided between the antenna terminal and transmitting filter, and between the antenna terminal and receiving filter, in order to heighten the impedance as seen from the antenna terminal in the rejection band, and thereby suppressing leak of a receiving signal toward the transmitting terminal and a transmitting signal toward the receiving terminal, respectively (that is, isolation is assured).\nIn this constitution, the transmitting signal entering from the transmitting terminal is output to the antenna terminal without leaking to the receiving terminal. Likewise, the receiving signal entering from the antenna terminal is input into the receiving terminal without leaking to the transmitting terminal.\nGenerally, as the transmitting filter, a band-rejection filter is used because it is enough to reject only the receiving band, while a band-pass filter is used for the receiving filter because it is required to reject not only the transmitting frequency, but also the neighborhood band of local oscillator frequency and image frequency of the receiving signal.\nTherefore, the ladder-type SAW filter of small size and high selectivity has optimum features as the receiving filter. However, when the SAW filter is used as the receiving filter, an electric power of about one watt is applied to the transmitting band which is the rejection band allocated in the lower outside of receiving band, and therefore a corresponding withstanding power is needed. As mentioned in detail below in the embodiment, when the circuit is composed in the serial-branch input type, the withstanding power to the transmitting power is enhanced, and problems are practically solved.\nHowever, when the conventional SAW filter 11 is used as the receiving filter for the antenna duplexer, since the reflection coefficient is small at the transmitting frequency to be rejected, the impedance when seeing the receiving filter from the antenna terminal through the phase shifter is not sufficiently high, and the isolation between transmitting terminal and receiving terminal may not be assured.\nThe invention is intended to solve the above problems, and it is hence an object thereof to realize not only a SAW filter capable of obtaining a high reflection coefficient in rejection band even in the case of the serial-branch input type with high withstanding power, but also a SAW filter capable of obtaining a high withstanding power characteristic against the transmitting power and assuring a sufficient isolation between the transmitting terminal and receiving terminal when used as the receiving filter for an antenna duplexer, and the object also includes to realize an antenna duplexer using the SAW filter and a mobile communication terminal using the same.\nTo solve the problems, the invention presents a ladder-type SAW filter composed by connecting alternately a serial branch having at least one SAW resonator connected in series to a signal path, and a parallel branch having at least one SAW resonator connected between the signal path and the ground, in which the sizes of wiring patterns of the parallel branch are smaller than wiring patterns of the serial branch.\nAccording to this constitution, even in the case of serial-branch input type, the inductive components of the parallel-branch wiring patterns are small, and hence a large reflection coefficient is obtained in the rejection band, thereby realizing a SAW filter with a high withstanding power characteristics and assuring a sufficient isolation between the transmitting terminal and receiving terminal, and also an antenna duplexer using the same and a mobile communication terminal using the same."}
{"text": "Cryogenic distillation is currently the preferred process for producing high purity oxygen (&gt;95%) in large scale plants (50-2000 ton/day). However, contaminants in the compressed air feed, i.e., water, carbon dioxide and trace hydrocarbons, must be removed before conducting the distillation process in order to prevent blocking of heat exchangers or distillation equipment and buildup of hazardous concentrations of hydrocarbons in the distillation column sump. Reversing heat exchangers are commonly employed to remove contaminants in the front end of the cryogenic plant wherein such contaminants are condensed in the exchanger passages and then removed with a waste gas stream. Alternately, adsorbent beds containing zeolites or alumina which require periodic regeneration are used to adsorb such contaminants. In addition, hydrocarbons must often be removed from the liquid oxygen sump by using an adsorbent such as silica gel. These methods lead to increased capital costs and inefficiencies in the overall separation processes.\nAlternate methods for recovering oxygen from an oxygen-containing gaseous mixture include vacuum swing adsorption (VSA) and pressure swing adsorption (PSA) processes which employ selective adsorption of various components instead of conventional cryogenic steps to separate the mixture. As in the case of cryogenic processes, one or more of carbon dioxide, water and hydrocarbons must be separated from the feedstock prior to running the process to avoid deleterious interaction with the adsorbents.\nTypical processes for removing carbon dioxide and water from oxygen-containing gaseous mixtures employ an adsorbent or desiccant and are capable of removing water vapor to very low levels, often to a dew point of less than -50.degree. F. These processes possess a drawback in that the adsorbent bed must be regenerated, usually by purging the adsorbent bed with a low pressure dry waste gas or by using some portion of the product stream if a suitable waste gas stream is not available. Consequently, these systems are operated in a cyclic manner requiring duplication of equipment, operation of automated, timed switching valves and separate heater devices. An unavoidable loss of the gaseous feed often occurs during regeneration of the adsorbent.\nU.S. Pat. No. 5,108,465 discloses a process for separating oxygen from an oxygen-containing gaseous mixture which comprises contacting the oxygen-containing gaseous mixture with a membrane which is impermeable to gas yet which is capable of conducting electrons and oxygen ions. The membranes are formed from a ceramic material selected from the group consisting of BaFe.sub.0.5 Co.sub.0.5 YO.sub.3 ; yellow lead oxide; ThO.sub.2 ; Sm.sub.2 O.sub.3 -doped ThO.sub.2 ; MoO.sub.3 -doped Bi.sub.2 O.sub.3 ; Er.sub.2 O.sub.3 -doped Bi.sub.2 O.sub.3 ; Gd.sub.2 Zr.sub.2 O.sub.7 ; CaTi.sub.1-x M.sub.x O.sub.3-.alpha. wherein M is Fe, Co or Ni, x is 0-0.5 and .alpha. is 0-0.5; SrCeO.sub.3 ; YBa.sub.2 Cu.sub.3 O.sub.7-.beta. wherein .beta. is 0-1 and (VO).sub.2 P.sub.2 O.sub.7.\nThe ceramic materials disclosed in U.S. Pat. No. 5,108,465 comprise ionically conductive materials, which are commonly used in fabricating solid oxide fuel cell components, and superconducting materials. For example, YBa.sub.2 Cu.sub.3 O.sub.7-.beta. is an ionically conductive superconducting material. However, barium-containing ceramic materials are well known to be adversely affected by the presence of carbon dioxide. For example, literature references teach that water and carbon dioxide will irreversibly degrade YBa.sub.2 Cu.sub.3 O.sub.7 destroying superconductivity properties of the material upon contact with carbon dioxide at temperatures greater than about 400.degree. C. This phenomena is discussed in numerous articles including those by E. A. Cooper et al., J. Mater. Res. 6 (1991) 1393 and Y. Gao et al., J. Mater. Res. 5 (1990) 1363. Like barium, other Group II metal oxides will react with carbon dioxide and water to form hydroxy carbonates. One of ordinary skill in the art would be lead to believe that use of membranes containing strontium, calcium or magnesium would be restricted to the separation of oxygen from oxygen-containing gaseous mixtures which are substantially free of carbon dioxide, water or hydrocarbons.\nConsiderable effort is being expended in developing an oxygen recovery process wherein the feedstock does not have to be pretreated to remove carbon dioxide, water or volatile hydrocarbons prior to conducting the separation. Moreover, improved processes are being sought to develop Group II-containing ceramics membranes which are capable of separating oxygen from oxygen-containing gaseous mixtures containing water or carbon dioxide."}
{"text": "The invention relates to a screen which is particularly intended for a non-invasive functional unit measuring the function of at least one organ through skin contact, the functional unit comprising a supporting structure, an electrical circuit for operation, and electrodes for the skin contact, the electrodes being operatively coupled to the electrical circuit and mechanically attached to the supporting structure.\nVital functions can be measured telemetrically using a non-invasive measuring device. An example of such a measuring device is, for example, a system which measures human heart rate and usually comprises as functional units a transmitter unit, a receiver unit and a data transfer unit. The transmitter unit refers to a transmitter unit equipped with electrodes and worn on the human body, particularly on the chest, the transmitter unit usually being implemented in the form of a transmitter belt attached around the body. The receiver unit refers for instance to a watch-like receiver unit which is worn on the wrist, the receiver unit having a telemetric or wired coupling with the transmitter unit. The data transfer unit transfers data received in the receiver unit to a computer, for example. The computer can also control the transmitter unit and the receiver unit via the data transfer unit.\nHeart rate measurement is based on monitoring the function of the heart. When the heart contracts, it generates a series of electric impulses that can be measured in the body. The measurement and analysis of this signal is called electrocardiography (ECG). The signal itself is called an ECG signal. Different phases of the heart cycle can be discerned in the ECG signal. These are called P, Q, R, S, T and U waves.\nThe transmitter part located next to the body suffers from electromagnetic interference, and particularly from problems caused by static electricity. A typical problem situation occurs when a user exercises with a shirt on. When the user moves, the shirt moves, flapping against the user\"\"s body and the transmitter unit. Synthetic fiber textiles in particular that are poor at conducting electricity attract electric charges specifically in dry weather, at least before the user starts sweating. In terms of electrical engineering, a great amount of charge at random amplitude and at random frequency then moves in the vicinity of the transmitter unit. Such a random movement of a great amount of charge is capacitively coupled to the transmitter unit, interfering with its operation.\nAn object of the invention is thus to provide a screen so as to solve the above problems caused by static electricity. This is achieved with a screen described in the introduction, the screen being characterized by being arranged to closely follow the electrode or the supporting structure at least at the electrode.\nA plurality of advantages can be achieved by the screen of the invention. The screen prevents disturbance from being connected to the functional unit; hence, the operation of the functional unit becomes more reliable."}
{"text": "It is generally known to those skilled in the art of golf-club making that a player's performance level may be enhanced by custom fitting a set of golf clubs to the player. These custom-fit golf clubs are selected, using anyone of a plurality of fitting methods, to improve the player's ability to consistently execute well-struck shots. Most fitting methods require a player to perform one or more golf swings using a single test club, e.g., a 6 iron. The collected data is then extrapolated to the rest of the clubs with the set. However, a player who is fitted into a custom set using a single test club may later discover a non-uniform shot-distance progression between the clubs within the set. Accordingly, the player may have difficulty selecting the appropriate golf clubs from the set for various shots during a golf round."}
{"text": "This application was originally filed as international application PCT/NL99/00498 on Aug. 4, 1999.\nThe present invention relates to a closing valve for a container or a pipe system, comprising a closing jacket connectable to an opening in the container or forming part of the container, which closing jacket is provided on the inside with a narrowed and a widened portion, and a valve part movable in this closing jacket with a closing element and a clamping element rigidly connected therewith and extending in the closing jacket, which clamping element is provided on at least the outside with a thickening which, when the passage through the closing jacket is sealed by the valve part, is in contact with the widened portion in the closing jacket, while, furthermore, a head part movable back and forth in the closing jacket, optionally to be put into circulation apart, is provided, by means of which the valve part can be moved with respect to the closing jacket such that the passage therethrough can be released and/or closed.\nSuch a container may be a cardboard, metal or plastic packing as well as a compressible plastic bag, a glass, bottle and the like. The container is suitable for gaseous products as well as all kinds of liquids, optionally mixed with gas, both under pressure and under vacuum, highly viscous substances, such as puree, and further granulates, granular material, etc. The closing valve may also be included in a pipe system.\nA closing valve as described in the opening paragraph is known from WO94/29215. In said document a thickening at the outside of the clamping element has the function to lock a head part with a tube relative to a valve part in the form of a closing plug. Because said thickening only temporarily rests against the edge of the thickening, for a short time a counterforce, necessary to lock the head part with the tube relative to the valve part, is exerted. By pushing the close jacket in the first instance said thickening meets the edge of the narrowed part of the closing jacket. At that moment the head part with the tube is still free removable (FIG. 2). When the head part with the tube is pushed further, the head part is locked in the xe2x80x98valvexe2x80x99 part because the edge of the valve part engages a groove between the head part and the tube (FIG. 3). When thereafter the head part with the tube is pushed further again, said thickening being pressed along the edge of the narrowed part of the closing jacket, the container is fully opened as then the outlet opening in the tube is located in the container (FIG. 4). In the latter position there is no locking between the valve part with the head part on the one hand and the closing jacket at the other hand, so that the position of these parts relative to each other is not determined any longer and not adjustable at all. As the valve element is in the form of a closing plug, the outlet opening is brought in the tube, with the risk that when the tube is wrenched off the valve part, the valve part will fall into the container, which then cannot be closed again.\nThe object of the invention is to provide for such a container an efficiently and inexpensively producible closing valve, via which the container can be easily filled and closed and the substances contained therein can be easily poured out or pressed or sucked out of the container, and by means of which closing valve the above problems are obviated.\nTo this end, according to the invention the closing valve is characterized in that, the thickening on the outside of the clamping element is located near the upper part of the clamping element, such that, when the passage through the closing jacket is released, the thickening on the outer side of the clamping element is brought into the narrowed portion of the closing jacket, as a result of which the clamping element reaches a position in which it is engaged by the head part.\nThe closing vale is therefore opened by pushing down the valve part by means of the head part, while by pulling up the head part, which as then engaged the clamping element, the valve part is taken along and the closing valve is gradually closed. The head part can therefore also be designed as a tap and brought into several positions and optionally fixed in these positions with respect to the closing jacket, so that gradually opening and closing the passage by the closing jacket becomes possible.\nAlthough it is sufficient for specific uses to only open the container by means of the head part, in which connection it is not important that the valve part lands in the container, it is important for many other uses that the head part is provided with locking means for preventing the valve part from being pushed completely out of the closing jacket.\nAs appears from the foregoing, the closing valve, in its basic structure, is composed of only three easily producible parts: a closing jacket, a valve part and a head part; moreover, an inexpensive final product can be obtained by, for instance, injection molding these parts from some hard plastic. Of course, it remains possible that the closing jacket forms part of the container and/or that the head part forms part of a draw-off tap.\nFor many uses the head part will be provided with a throughbore which, at the lower end, is in open communication with the inside of the closing jacket, so that, when the closing valve is opened by pushing in the head part, an open communication with the atmosphere results, to which optionally other filling, outlet or pouring means or the like can be connected. In a favorable embodiment the head part is provided with a widened upper part which, when the valve part is pushed out of the closing jacket at least partially, seals the upper side of the closing element.\nIn addition to the three above basic parts, further parts may be provided, such as a sealing element to be arranged between the closing jacket and the head part, in particular an O-ring, to enable a fluid-tight closure. Further sealing elements, for instance around the peripheral edge of the closing element, may be provided; these sealing elements, however, may also be formed by, for instance, plastic edges forming part of the relevant components.\nAlthough the clamping element may be cylindrical, while an annular thickening may be provided on the outside as well as on the inside, it is difficult to push in such a clamping element on all sides; anyhow, the clamping element then needs to be rather easily deformable. It is better when the clamping element is finger-shaped, with each of the finger-shaped parts on the inside and outside being provided with a thickening. It is indeed easier to push in the individual fingers. The thickening on the inside then serves to facilitate engagement by the head part.\nWhen the container is formed by an easily compressible packing, such as, for instance, a plastic bag to which the closing valve as hitherto defined is connected, and furthermore, this easily compressible packing is arranged in an outer packing which is not compressible or is relatively difficult to compress, the closing valve needs to be provided with additional means. To this end, according to the invention a collar element connectable to an outer envelope or forming part thereof is provided, which collar element fits around the closing jacket, while between the outer circumferential surface of the closing jacket and the inner circumferential surface of the collar element gas supply channels, in particular air supply channels, are provided, which are opened when releasing the passage in the closing jacket. This has the advantage that gas can flow into the space between the inner packing, this is the easily compressible packing, and the outer packing, so that, when a pump is connected to the head part, the inner packing can be easily emptied by sucking, while the inner packing is compressed. Also, gas can be forcibly brought into the above space to enable the contents to be pressed out of the inner packing. This is important when the head part is formed by a filling head or a tapping head of, for instance, a tapping plant for beverages."}
{"text": "(a) Field\nExemplary embodiments of the invention relate to a quantum rod sheet, a backlight unit including the quantum rod sheet, a display device including the backlight unit and a manufacturing method thereof.\n(b) Description of the Related Art\nVarious types of flat panel display, such as a liquid crystal display (“LCD”), a plasma display panel (“PDP”), an organic light emitting display (“OLED”), an electrowetting display (“EWD”), an electrophoretic display (“EPD”), an embedded micro-cavity display (“EMD”), and a nano-crystal display (“NCD)”, have been developed from a cathode ray tube (“CRT”) type using a cathode ray tube.\nAmong the various types of flat panel display, the LCD, which is one of the most widely used types, has characteristics, such as miniaturization, weight lightening, and low power consumption. In general, the LCD is a device, in which an electric field is generated in a liquid crystal material disposed between an upper substrate and a lower substrate by applying different potentials to a pixel electrode and a common electrode thereof such that an arrangement of liquid crystal molecules is changed, and transmittance of light is thereby controlled to display an image.\nIn the LCD, a liquid crystal panel does not emit light therefrom, and the liquid crystal panel thereby includes a backlight unit for providing light to the liquid crystal panel. Among the various types of flat panel display, other light receiving type display devices may include the backlight unit.\nThe backlight unit typically includes a light source and a plurality of optical sheets to improve luminance of light provided from the light source and to distribute the light substantially uniformly over an entire region thereof. Among the optical sheets, an optical sheet that may be used to improve the luminance (e.g., a luminance enhancement film) is typically substantially expensive, which may lead to an increase in manufacturing costs of the display device.\nIn a backlight unit configured to generate white light, the backlight unit typically includes a light emitting diode (“LED”) that emits light having a specific wavelength and a phosphor that changes the light having the specific wavelength emitted from the LED into the white light to provide the white light."}
{"text": "Occasions arise where a well has to be plugged and abandoned. Typically, such wells will have a cemented casing and another tubular string inside the casing. The plug and abandon procedure involves severing the inner tubular followed by axially milling the severed tubular for a predetermined axial distance followed by setting a plug and depositing cement on top of the plug. The tubular to be severed has a nominal outer dimension between upsets at opposed ends where the joints occur. The axial milling has to cut the nominal dimension as well as the upset outer dimension. In either case it is not desirable to cut into the surrounding tubular as the integrity of the surrounding tubular is important to the functioning of the plug created as part of the abandonment process. On some occasions the tubular to be severed and axially milled is not disposed concentrically in the surrounding casing and sometimes the inner tubular lies on the casing. Axially milling the severed tubular involves using blades that extend a predetermined amount that is short of reaching the surrounding tubular.\nThe way this process is done today takes multiple trips. The first trip is with a mill with extending blades whose pointed ends radially penetrate the wall of a tubular until that wall is breached. These blades are limited in their radial extent to avoid damage to the surrounding tubular. The tool is then tripped out for a blade replacement to dress the tool with blades that will axially mill the nominal outside diameter of the tubular in a location where the severing took place, which is generally between upsets that have an even larger outer dimension. The axial milling continues between the upsets until an upset is reached. At that point there is another trip out of the hole to redress the tool with longer axial milling bladed that can reach further to the exterior of the larger upset dimension and yet not far enough to gouge the surrounding tubular. Should further axial milling be needed after an upset then another trip to swap back to the shorter axial milling blades is needed.\nIn one embodiment of the invention all three blade types are provided on a tool housing run into the tubular to be severed and then axially milled. The severing blades are initially extended for the radial cutting through the wall after which those blades are no longer needed. With the aid of specially configured darts that land in a profile either the shorter or the longer axial milling blades are extended. The uniquely configured darts are blown out into a catcher to allow selection of the longer or the shorter blades alternatively as many times as needed depending on the length of the axial milling distance and how many upsets need milling in that desired interval. No matter if the shorter or longer axial milling blades are extended, there is an open circulation path for the extended blades to cool them and for cuttings removal to a collection device or to the surface. While severing and milling is the preferred application, other uses are envisioned where there is a need to sequentially operate a variety of tools in a single trip particularly when two of the tools need to be alternatively operated multiple times. While the preferred actuation system is described in a hydraulic operating environment other ways to achieve the desired tool operating sequence are also contemplated with the objective being an economical and hence reliable design that can perform the required actions in a single trip and preferably without well intervention.\nU.S. Pat. No. 8,955,597 shows a sequential use of a mill to make a radial cut with blades 76 followed by extension of a reamer 84 to enlarge the opening where the casing was earlier milled away. This application is limited to a single string in the hole and the sequence is executed but a single time. U.S. Pat. No. 5,765,640 describes a multipurpose tool where one or more pistons can be used to operate multiple tools at the same or different times but shows or describes no structure as to how that can be done one time, not to mention multiple times. U.S. Pat. No. 8,141,627 shows multiple rows of cleaning blades that extend at the same time and are spring retracted.\nThose skilled in the art will have a better understanding of the preferred embodiment of the invention from the description below and the associated drawings while recognizing that the full scope of the invention is to be determined from the appended claims."}
{"text": "Compared with traditional lighting technologies such as tungsten bulbs, fluorescent tubes and LEDs, OLEDs have the advantages of having slim physical sizes, wide color variety, and physical bendability. These characteristics of OLEDs avail the technology to many new applications, such as decorative and effect lighting, which were not practically achievable in the past. For example, an OLED lighting module, without the accompanying power convertor and mechanical supporting structure, may have a thickness of less than 1 mm. On the other hand, its relatively high price at the present limits its application in general lighting purposes.\nSimilar to LED lighting and unlike tungsten bulbs and fluorescent tubes, OLED lighting requires to be driven by electric current as power source and brightness control. Therefore, a power converter is required to provide a controlled electric current. In the applications of decorative OLED lighting, the power sources are often 12-48V DC voltage sources. A power converter is then used to convert the constant voltage to a controlled constant current. A conventional power converter uses a single-phase buck current converter as shown in FIG. 1. It converts a direct current (DC) voltage to the controlled constant current.\nTo maintain the slimness and bendability of an OLED lighting panel (i.e. within 1.5 mm thickness including mechanical supporting structure), the power converter needs to be thin as well. Otherwise, the power module cannot be integrated into the OLED panel and must be installed separately and connected to the OLED panel with long wires. Long wires carrying constant current are not desirable as there can be substantial energy loss in the transmission. More importantly, the wires of individual panels have to be separated due to its constant current nature, and cannot be combined to save space and cost.\nTo build a 1.5 mm thick power module, all components needed to be 1.0 mm or less in height as 0.5 mm is normally reserved for the flexible printed circuit board (PCB) and mechanical mount. Among the electronic components within a power module, inductors with low profiles are the least available components. For instance, a normal surface mount 15 uH 0.6 A inductor already has a thickness of 1 mm. As another example, a low-profile 60 uH 1.2 A inductor has a thickness of 5 mm, while an inductor with a normal profile is even thicker.\nOn the other hand, an OLED lighting panel usually has a large surface area (i.e. 100 mm×100 mm is commonly found and 300 mm×100 mm is forthcoming) that can easily conceal a power module there behind. This means that there are more space allowance in the length and width dimensions, but very limited in the height dimension for the power module. As such, the inductor's thickness is the critical factor, and there is no practical solution in the art currently to enable slim power modules for OLED lighting due critically to the thickness of normal inductors.\nA possible but undesirable solution is to replace one thick inductor by multiple smaller thin inductors, so that the thickness of the power converter can be reduced in the expense of having a larger occupying area in the length and width dimensions. For example, as shown in FIG. 2, the 60 uH 1.2 A inductor is replaced by a network of sixteen 15 uH 0.6 A inductors. Since the height of each 15 uH 0.6 A inductor is only about 1.0 mm, the power converter with this network of inductors may meet the thickness requirement. However, the total magnetic material used is doubled as compared with the single 60 uH 1.2 A inductor. Also, the much higher component count increases material cost and assembly cost significantly. The amount of magnetic material used is a major factor affecting the size of the inductors."}
{"text": "An underground mining operation such as a coal mine presents hazardous conditions for both machines and workers. An underground coal mine resembles a large construction project and frequently groups of workers are around or near high power electrical equipment. Mining equipment quite often requires large amounts of electrical power located near the underground excavation point which creates a need for large power transformers near the excavation point. These electrical power transformers combined with input, output, and control circuitry are combined in one enclosure called a power center. A power center is typically protected by a solid sheet-metal enclosure. The enclosure's main purpose is to protect the electrical components from dust, dripping water, and condensation, and dissipate heat. In addition, the enclosure must protect the electrical components from damage due to rock slides, rock falls, and cave-ins.\nIn addition to protecting the internal electrical components, enclosures are designed to increase mine safety. In the event of an internal electrical failure within the power center, the enclosure must retard the introduction of flames or sparks into the mining area where flammable materials may be ignited.\nCurrently, arc resistant enclosures in the industrial markets are designed to withstand the pressures and temperatures of gasses associated with an internal arcing fault and are configured to safely direct these gasses out of the roof of the enclosure away from personnel. This type of configuration is problematic for the underground mining industry where the ceiling of the mine is often only a few inches away from the top of the electrical enclosure. To make things worse, the ceiling and walls of the mine often contain flammable material such as coal. The surrounding air is also often times contaminated with coal dust and possibly explosive gases."}
{"text": "In various situations there is a need to distribute heat in outdoor areas. For example, in horticultural applications crop yields can be devastated by frost damage experienced at the wrong time of a growing cycle. Heavy rainfall and water sitting on the skin of a fruit can also impact on the quality of the produce generated by a grower. Crop yields are also dependent on successful pollination, with ambient temperatures having a role to play in this process with various crops.\nSome attempts have been made to control the temperature experienced inside orchards and other similar crops. One technique currently used is the distribution of large numbers of localised heating devices or pots over the orchard floor. However, these stationery heat sources only distribute heat very locally and with a limited lateral extent. This approach requires a large number of individual heat sources to be laid out and activated over the extent of the orchard, and then subsequently collected again after the weather improves.\nAn alternative approach is to tow a heat source connected to a fan through an orchard. Although this approach does not require large number of individual heat pots to be deployed and collected, it does require a driver to continuously run a vehicle over and through the orchard—potentially through the early hours of the morning when frost risk is at its greatest. Furthermore, these towed heating systems again have a limited ability to distribute warm air laterally and generally only provide a narrow swathe through the orchard with little persistence. If the heat is too great close to the plants it can damage the crop. With lack persistence the heater needs to operate over a repeating track that returns within about four minutes which severely limits the area that a machine can cover effectively.\nIn some circumstances a temperature inversion may occur in the air surrounding an orchard. In these circumstances an upper warm air layer forms above the orchard while colder air is resident in and on the plants. Helicopters and wind turbines had previously been used to force the warm air above the orchard down into the plants, and although this approach is effective in preventing frost and generally raising orchard air temperatures, is also expensive in terms the cost of machinery required and the energy consumed by this machinery. Furthermore, this approach generates a significant amount of noise at time periods where local residents may be sleeping.\nIt would therefore be of advantage to have an improved heat distribution apparatus and method of operating same which addressed any or all of the above issues. A heat distribution apparatus and method of operating same which provided the public with at least an alternative choice would also be of advantage. Improvements over the prior art which allow for the distribution of warm air in or above the canopy of an orchard would also be of advantage, as would improvements which allowed for control of the volume and temperature of air distributed to various sections of an orchard."}
{"text": "Protection gear for various purposes has been used throughout history. Hunting, traveling, warfare, and daily occupations led to the development of leather or fabric guards designed to protect parts of the body exposed to danger. These early body shields evolved as needs arose and materials developed, until today specialized protective gear exists for almost every imaginable human activity.\nOf particular interest is the use of protective devices for the knee and lower leg. Such devices are typically used in skilled construction trades, for example, in the installation of tile and other floor coverings. These professions require their practitioners to spend long, strenuous hours of kneeling and crawling, and knee problems are considered to be an almost inevitable occupational hazard.\nKnee protection devices for these and other purposes have been the subject of a high degree of inventive activity, as evidenced by the patent literature. For example, U.S. Pat. No. 7,451,493 shows a knee pad construction including a shaped cushion element with a concave interior with a recess having a flowable gel insert therein. The recess and insert are shaped to accommodate the left knee or the right knee of an individual.\nU.S. Pat. No. 7,448,088 deals with a knee pad for attachment to a person's knee cap and an adjacent portion of the person's lower leg, comprising a cushioning pad that includes a forward section and a rearward section. The forward section has a knee well for receiving the person's knee cap and the knee well is off-center in the forward section. The rearward section of the cushioning pad has an interior concavity for receiving the person's lower leg.\nU.S. Pat. No. 7,237,270 shows a garment having a pants leg with a front layer stitched along an outer peripheral seam to the pants leg to form a pocket with an interior cavity. The front layer has a central opening spaced from the outer peripheral seam and defined by an inner periphery. A protective insert, which may be formed of foam, and which is larger than the central opening, is removably positionable within the pocket. A stiff cap formed of a material such as SANTOPRENE® thermoplastic elastomer, is stitched to the insert. The cap has a flange which projects outwardly from a central region around a line of attachment. The central region is no larger than the pocket central opening. The insert is removably receivable within the pocket such that the inner periphery of the front layer is engaged between the stiff cap flange and the insert.\nU.S. Pat. No. 7,197,770 involves a rolling kneepad having a weight transfer support assembly with swivel casters to support a user, a kneepad cushion provided above the support assembly to comfortably receive a user's knee, and a shin guard peripherally attached to the support assembly and extending outwardly therefrom. The elongated shin guard has a distal end which abuts the upper surface of a user's shoe or boot. Consequently, when the user moves from a kneeling to a standing position, the weight of the rolling kneepad is supported on the top surface of the user's shoe or boot. Downward slippage of the rolling kneepad along the user's leg is prevented and manual reposition of the rolling kneepad prior to the user resuming a kneeling position is avoided. Straps having hook and loop fasteners at the ends are typically used to affix the rolling kneepad to the body of the user.\nU.S. Pat. No. 7,096,507 describes a lower leg appliance providing padding and cushion support for the knee, ankle and hip joints in the leg, by providing a hard surfaced knee pad, a hard surfaced ankle pad and a soft surface calf and thigh pad positioned between the calf and thigh, held in place by at least three adjustable straps having a closure means, the three pads working in conjunction to provide relief and support to a persons knees, ankles and hips while squatting or kneeling during chores, activities or labor.\nU.S. Pat. No. 6,795,974 is directed to a kneeling pad assembly having an articulated supporting spine assembly with an upper spine member carrying a protected cushioned kneecup and a lower spine member carrying a shin engaging cushion. The pad assembly is supported on the leg of a user only by a pair of leg straps extending from the lower spine member behind the user's leg below the knee. The upper spine member is cantilevered above the lower spine and is constructed to be pivotally moveable with a snap action between two stable positions. The first stable position is with the upper spine collinear with the lower spine to hold the cushioned kneecup against the user's knee. The second stable position is with the upper spine member angled away from the user's leg at an acute angle to the lower spine member to hold the kneecap out of contact with the user's knee to avoid discomfort and displacement of the kneeling pad assembly during standing or walking.\nU.S. Pat. No. 6,178,555 discloses a lower leg and foot cover for covering the knee and lower leg regions of a user's pants and the front of the user's shoes from soiling and soaking. The lower leg and foot cover includes a knee cover which is coupled to an upper end of a flexible lower leg cover. A flexible shoe cover is coupled to a lower end of the lower leg cover. Flexible straps are provided for securing the knee cover and the lower leg cover to the wearer's leg.\nU.S. Pat. No. 6,405,383 is directed to a leg protective system selectively covering a person's lower leg and knee for protective against environmental elements. The leg protective system comprises a first protective cover and a second protective cover which are selectively attachable to each other to define a single leg protective system. The first protective cover is adapted to cover the front and side portions of a wearer's knee. The first protective cover has a top portion and a bottom portion, whereby the top and bottom portions are formed of a flexible material for generally fitting the contour of the knee. The bottom portion has a first releasable strap assembly adapted to extend behind the wearer's leg to operatively attach the first protective cover over the knee. The second protection cover is adapted to cover the front and side portions of a wearer's lower leg. The second protective cover is formed of a flexible material for generally fitting the contour of the lower leg. The second protective cover has a second releasable strap assembly adapted to extend behind the wearer's lower leg to operatively attach the second protective cover to the lower leg. The first and second protection covers are selectively attachable adjacent their respective bottom and top edges to define, when so attached, a single leg system.\nU.S. Pat. No. 5,652,956 discusses a shin pad in which the shin shield is in two pieces, namely an upper piece which is secured to the lining of the knee cap, or to the knee cap itself, and a lower piece which is slidable up and down relative to the upper piece, and to which a lower liner is secured, the lower liner being independent of the knee liner. Preferably, the lower piece is slidable with respect to the upper piece by virtue of one or more slotted holes in either or both of the upper piece and/or lower piece, with bolts and T-nuts being used in the slots to secure the pieces together in the desired position. Preferably, there are three such slotted holes, namely a front central one and two side ones, namely one towards either side of the pad. Preferably, there is sufficient sloppiness in the slots to provide a varus/valgus adjustment, i.e. an adjustment of the angle of the shin shield, in the vertical fore-and-aft plane, relative to the orientation of the knee cap, by permitting the lower piece to be secured with one side higher than the other relative to the upper piece.\nU.S. Pat. No. 5,090,055 concerns a unique protective knee pad which is comprised of one-piece molded resilient polyurethane foam. The pad has a special U-shaped sole with a transverse bar through the middle portion and a transverse bar at the rear. Two cavities are created by the transverse bars which are adapted to entrap air when the wearer is in the operative kneeling position thereby cushioning the delicate parts of the patella and upper shin, by preventing the outwardly extending portion of the knee from coming into contact with the ground.\nU.S. Patent Publication No. 20080168589 deals with a knee pad with a support frame, a pad, and a boot. The pad is assembled on the support frame to provide a cushioning surface between the support frame and the wearer's leg. The pad has ears in the knee section and in the ankle section, which serve to maintain the support frame and pad in proper alignment on the leg. The boot provides a non-slip, non-marring surface against the floor, when the wearer is working in a kneeling position. A moisture-absorbent liner, which is easily removed for cleaning or replacement, may also be fitted on top of the pad, to provide an absorbent layer between the pad and the wearer's leg.\nU.S. Patent Publication No. 20080072359 shows protective knee pads. The knee pads are comprised of a first shell member which fits over the patella and knee joint and a second member which is attached to the lower thigh and pivotally engaged with the knee pad shell so that they may articulate one with respect to the other while providing additional support and protection for the user of the knee pad. An additional shin pad may also be attached to the lower end of the shell which covers the knee.\nU.S. Patent Publication No. 20070050877 shows a clothing adherable knee pad for a pair of shorts including a bendable flexible arcuate pad which bends over the knee of a kneeling person. The knee pad is cantilevered downward from at least two linearly extending elastic, stretchable straps which are connected at distal ends thereof each to a respective clothing engageable clip at upper region of the knee pad, wherein each clip is attachable to a leg of a pair of short pants. The two or more linearly extending elastic, stretchable straps engage the knee pad by being attached thereto or by being woven through respective pairs of slits extending through the upper region of the knee pad. The straps are of sufficient length so that when attached to clothing, such as short pants, respective locations of the clips upon short pants legs avoids uncomfortable encroachments to respective lateral and rear areas of the knee of the user.\nIt can be seen from the foregoing that the need exists for knee pad that provides a durable, inexpensive, easily manufactured, and comfortable pad adapted to a wide variety of potential uses."}
{"text": "Wireless communication systems that use pilot information are well known in the art. Pilot information generally serves to facilitate an assessment of the wireless channel by a receiving party. By providing pilot information having known characteristics, a receiver can utilize its knowledge of the differences (such as amplitude and phase differences) between the pilot information as it should have been received and the pilot information as it was received to ascertain and compensate for alterations that were occasioned by the channel to thereby promote a more accurate recovery of other information as was also received by the receiver via that channel. Pilot information can also be used to effect other purposes as well.\nUnfortunately, there are a variety of trade-offs that are ordinarily accepted when combining pilot information with a user payload. Such trade-offs, in various ways, usually impact the usable bandwidth or user payload throughput capacity of the channel. For example, in a code division multiplexed wireless communication system that makes use of orthogonal variable spreading factor codes to permit differentiation of a number of data channels, a control channel that includes pilot information will also have, or share, such a code. Within the context of such a system, this means that some portion of overall available power must be allocated for the control channel/pilot information. When less power than is necessary at any given moment is used to support the pilot information, however, user information throughput capacity is negatively impacted.\nConversely, of course, the opposite can also occur. More power may be accorded the pilot information than is necessary to overcome the present vagaries of the channel. As a result, less power is available for the user payload and this again can lead to a diminution of the receiver's ability to properly decode the user payload.\nOne prior art approach relies upon a simulated study of a channel for a given set of conditions (such as, for example, a given data rate and service overhead for the data payload). Running the simulation develops a set of corresponding gain factors to be used with the data/control channel. By providing these results in a look-up table, the gain factor as corresponds to a given data rate and service can then be recovered and used during ordinary operations. One significant drawback of this approach is that this technique does nothing to assist a wireless communication unit that must use, for whatever reason, a data rate that is not represented in the look-up table. Since new data rates are a common and frequent occurrence, this constitutes a considerable obstacle to effective use.\nSkilled artisans will appreciate that elements in the figures are illustrated for simplicity and clarity and have not necessarily been drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help to improve understanding of various embodiments of the present invention. Also, common but well-understood elements that are useful or necessary in a commercially feasible embodiment are typically not depicted in order to facilitate a less obstructed view of these various embodiments of the present invention."}
{"text": "Fuel cell stacks are electrochemical devices that produce water and an electric potential from a fuel, which typically is a proton-liberating source, and an oxidant. Many conventional fuel cell stacks utilize hydrogen gas as the proton source and oxygen, air, or oxygen-enriched air as the oxidant. Fuel cell stacks typically include many fuels cells that are fluidly and electrically coupled together, often between common end plates. The fuel cell stack receives flows of hydrogen gas and air from suitable sources and distributes these flows to the individual fuel cells in the stack. The fuel cell stack includes manifolds and other delivery conduits to deliver and remove fluids to and from the fuel cells within the fuel cell stack. Conventionally, a fuel cell stack includes current collectors that are adapted to be electrically connected to an external load so that the electrical output produced by the fuel cell stack may be used to satisfy the applied external load.\nThe fuel cells in the fuel cell stack include anode and cathode regions that are separated by an electrolytic barrier, which may take the form of an electrolytic membrane. Hydrogen gas is delivered to the anode region, and oxygen gas is delivered to the cathode region. Protons from the hydrogen gas are drawn through the electrolytic membrane to the anode region, where water is formed. Conventionally, the anode and cathode regions are periodically or continuously purged to remove water and accumulated gases in these regions. While protons may pass through the membranes, electrons cannot. Instead, the electrons that are liberated by the passing of the protons through the membranes travel through an external circuit to form an electric current.\nA factor that affects the performance, or efficiency, of a fuel cell stack to produce an electric current is the temperature of the stack. Accordingly, it is desirable to maintain the fuel cell stack within a range of suitable operating temperatures, such as within upper and lower threshold temperatures. During startup, the stack may not be at a temperature within this desired temperature range, and in such a situation, it is desirable to transition the stack to a temperature within this range. Maintaining and/or starting up a fuel cell stack in subfreezing temperatures is more difficult because the liquid heat exchange fluid may be frozen or freeze when drawn from a reservoir for delivery to the fuel cell stack."}
{"text": "This disclosure relates to refrigeration systems for perishable items."}
{"text": "In recent years, biosensors have been contemplated as analytical alternatives to conventional methods in different fields. A biosensor is an analysis device made up of two fundamental elements: a bioreceptor (an antibody, a DNA probe, or a cell . . . ) prepared to specifically detect a substance, taking advantage of the specificity of biomolecular interactions, and a transducer or sensor, capable of interpreting the biological recognition reaction produced by the receptor and “transducing it” into a quantifiable optical or electrical signal. The most outstanding features of these devices which make them highly attractive options as analytical tools are their specificity, high sensitivity, response capacity leading to a short analysis time, their capacity of being included in integrated systems, ease of automation, capacity of working in real time, their versatility and low cost.\nProgress in the field of biosensors lies in the experience acquired throughout the years concerning the recognition capacity and properties of various biomolecules. Many biological instruments ranging from the simplest ones such as enzymes or antibodies to more complex genetic engineering products have been used as recognition elements. On the other hand, the recent advances in microelectronics, nanotechnology and the unique properties of specific materials have been key for such devices.\nDespite the foregoing, the current sensing methods are not always capable of meeting the requirements for reliability and speed. The time necessary to perform the assay and the sensitivity of the technique are some of the most significant limitations. Although the development of these devices has mainly been focused on the field of clinical diagnosis, their interest in other fields of application including the environmental, agro-food, chemical, pharmaceutical and military fields is on the rise today.\nThe present invention proposes a new biosensor based on the light to heat conversion properties of metal nanoparticles."}
{"text": "1. Technical Field\nThis invention relates to gas turbine engines in general, and to methods and apparatus for cooling a substrate exposed to high temperature gas in particular.\n2. Background Information\nEfficiency is a primary concern in the design of any gas turbine engine. Historically, one of the principle techniques for increasing efficiency has been to increase the core gas path temperatures within the engine. Core gas refers to air worked within the compressor that is mixed with fuel and combusted within the combustor. The increased gas path temperatures have been accommodated by using internally cooled components made from high temperature capacity alloys. Turbine stator vanes and blades, for example, are typically cooled using compressor air worked to a higher pressure, but still at a lower temperature than that of the core gas flow passing by the blade or vane. The higher pressure provides the energy necessary to push the air through the component. A significant percentage of the work imparted to the air bled from the compressor, however, is lost during the cooling process. The lost work does not add to the thrust of the engine and therefore negatively effects the overall efficiency of the engine. A person of skill in the art will recognize, therefore, that there is a tension between the efficiency gained from higher core gas path temperatures and the concomitant need to cool turbine components and the efficiency lost from bleeding air to perform that cooling.\nThere is, accordingly, great value in maximizing the cooling effectiveness of whatever cooling air is used. Prior art coolable airfoils typically include a plurality of internal cavities, which are supplied with cooling air. The cooling air passes through the wall of the airfoil (or the platform) and transfers thermal energy away from the airfoil in the process. The manner in which the cooling air passes through the airfoil wall is critical to the efficiency of the process. In some instances, cooling air is passed through straight or diffused cooling apertures to convectively cool the wall and establish an external film of cooling air. A minimal pressure drop is typically required across these type cooling apertures to minimize the amount of cooling air that is immediately lost to the free-stream hot core gas passing by the airfoil. The minimal pressure drop is usually produced through a plurality of cavities within the airfoil connected by a plurality of metering holes. Too small a pressure drop across the airfoil wall can result in undesirable hot core gas in-flow. In all cases, the minimal dwell time in the cooling aperture as well as the size of the cooling aperture make this type of convective cooling relatively inefficient.\nSome airfoils convectively cool by passing cooling air through passages disposed within a wall or platform. Typically, those passages extend a significant distance within the wall or platform along a substantially straight line. There are several potential problems with this type of cooling scheme. First, the heat transfer rate between the passage walls and the cooling air decreases markedly as a function of distance traveled within the passage. As a result, cooling air flow adequately cooling the beginning of the passage may not adequately cool the end of the passage. If the cooling air flow is increased to provide adequate cooling at the end of the passage, the beginning of the passage may be excessively cooled, consequently wasting cooling air. Second, the thermal profile of an airfoil is typically non-uniform and will contain regions exposed to a greater or lesser thermal load. The prior art internal cooling passages extending a significant distance within an airfoil wall or a platform typically span one or more regions having disparate thermal loads. Similar to the situation described above, providing a cooling flow adequate to cool the region with the greatest thermal load can result in other regions along the passage being excessively cooled.\nWhat is needed, therefore, is a method and apparatus for cooling a substrate within gas turbine engine that adequately cools the substrate using a minimal amount of cooling air and one that provides heat transfer where it is needed."}
{"text": "The present disclosure relates generally computer vision systems and relates more particularly to sensors for measuring the distance to an object or point in space.\nUnmanned vehicles, such as robotic vehicles and drones, typically rely on computer vision systems for obstacle detection and navigation in the surrounding environment. These computer vision systems, in turn, typically rely on various sensors that acquire visual data from the surrounding environment, which the computer vision systems process in order to gather information about the surrounding environment. For instance, data acquired via one or more imaging sensors may be used to determine the distance from the vehicle to a particular object or point in the surrounding environment."}
{"text": "The rising volume of communication, as a result of an increasing number of communications subscribers and rising demands on the volume of data to be transmitted, places greater and greater demands on switching centers, in particular on private branch exchanges, with regard to the volume of data to be transmitted for each communications link and the number of communications terminals which are to be connected to one another.\nCurrent devices are based, by way of example, on the TDM method (Time Division Multiplexing), in which communications data on different connections are transmitted in respectively defined timeslots. A connection between different communication parties is set up by a switching matrix which assigns incoming timeslots on an incoming connection outgoing timeslots to an outgoing connection on the basis of a control information item. Such switching matrices are generally of fixed proportions and can set up only a defined number of connections, which often makes it difficult to match exchanges to requirements. Another problem with such devices is that the timeslots are able to hold only a limited amount of data.\nOn account of different strengths and weaknesses of networks when transporting voice and data in the local domain, various communications networks have become established for specific purposes of use.\nFIG. 1 shows an example of a known private branch exchange 150 with two peripheral devices P1 and P2 to which a communications terminal KE1 and KE2 operating on a digital or analog basis is respectively connected. These peripheral devices P1 and P2 are accommodated in the same physical area as the central device ZE1. By way of example, they are in the same room or in the same cabinet as it. The terminals occupy defined timeslots in the PCM data stream (Pulse Code Modulation) with communication data. The digital or analog communications terminals KE1 and KE2 are respectively connected to subscriber line modules SLMO1 and SLMO2 which supply or take digital data, intended for the respective terminals or coming from the respective terminals, to/from the PCM data stream using timeslots stipulated by signaling. These PCM data streams are denoted by 100 and 200 in FIG. 1. In addition, signaling connections are shown which are represented by 110 and 210. It should be noted that these involve a logical representation and not a physical representation. In reality, however, the transport data and the signaling data are transmitted in the same connecting cable.\nThis figure also shows peripheral devices P1 and P2 and supply modules LTUC1 and LTUC2 which regulate the data traffic to the subscriber line modules of the respective peripheral devices. The peripheral device P1 is supplied with signaling data via the line 110, and the peripheral device P2 is supplied with signaling data via the signaling line 210.\nAs can clearly be seen here, this arrangement involves both the information to be transported and the signaling information being supplied to a central device ZE1. In this context, a messaging device DCL collects and distributes messages 2 which are interchanged between the central device ZE1 and the peripheral devices P1, P2. The call processing section CP controls the setup and clear-down of connections and, to this end, uses equipment-specific interface functions DH, inter alia, which are in the form of program modules, for example. This involves producing setting instructions 1 for the switching matrix MTS. These setting instructions essentially indicate which input of the switching matrix is to be connected to which output in order to provide a communications link. Control and connection functions are thus performed by a single physically incorporated functional unit in the communications network.\nProblems arise with such configurations because the data to be transported need to be supplied to the central device ZE1. This is the case even if, by way of example, two communications terminals which are connected to the same peripheral device P1 need to communicate with one another. The wiring complexity required for such devices increases with increasing distance between the terminals and the central device ZE1, which means that this type of arrangement restricts the extent of a private branch exchange or makes installation much more expensive when covering relatively large areas.\nIn such devices, problems likewise arise with regard to modular extendibility both in terms of the number of connections and in terms of the volume of data to be transmitted. This type of embodiment does not allow different data rates for each individual communications link.\nAlignment and/or transparent integration of different network infrastructures is likewise not possible."}
{"text": "An all-solid-state lithium secondary battery in which a solid electrolyte (e.g., a sulfide- or oxide-based solid electrolyte) is used as an electrolyte, has many advantageous such as high energy density and high power output.\nFor such an all-solid-state lithium secondary battery, it is known that active materials used in the electrodes are likely to be unstable and deteriorate. To avoid the problem, an electrode active material in which the surface is coated and protected with another component, is used. However, an all-solid-state lithium secondary battery using an electrode active material in which the surface is coated with lithium niobate, poses a new problem in that there is an increase in the reaction resistance.\nIn Patent Literature 1, a method for producing a cathode active material (a lithium-transition metal oxide powder) is disclosed, in which the surface is coated with lithium niobate by use of an aqueous solution containing a niobic acid complex and a lithium compound. According to the production method in Patent Literature 1, it is described that the content of carbon contained in the lithium-transition metal oxide powder in which the surface is coated with the lithium niobate, can be controlled to 0.03% by mass or less, so that carbon-induced lithium ion transfer inhibition can be suppressed, and the electrical resistance value (powder compact resistance) of the cathode active material can be decreased.\nPatent Literature 1: Japanese Patent Application Laid-Open No. 2012-74240"}
{"text": "High voltage equipment is commonly sealed within a welded enclosure. The seal is required because the enclosure contains either a gaseous (e.g., SF.sub.6) or liquid (e.g., oil) dielectric atmosphere within which the equipment operates. To enable connection between the enclosed equipment and external power lines, a bushing is used to provide a conductor through the wall of the enclosure while simultaneously insulating the conductor from the enclosure and maintaining the seal.\nBushings typically consist of a unitary metal conductor encapsulated in epoxy insulation. The inside end of the bushing (i.e., the proximal end of the conductor) is usually connected rigidly (e.g., by a metallic clamp and hardware) to the enclosed high voltage equipment, while the outside end of the bushing is designed to matingly receive a connector device, such as an elbow connector. The outside end of the bushing (i.e., the insulator body and/or the bushing-to-elbow interface) is subject to damage during both transportation and use which quite often renders the bushing unusable. In this event, the bushing must be replaced, which for a unitary bushing requires opening of the sealed enclosure and disconnection of the bushing from the high voltage equipment. Such replacement procedures are costly, time-consuming and often quite difficult to perform, especially when the equipment is located in the close quarters of an underground vault."}
{"text": "It has been common practice to include in lubricant formulations additives to provide improved antiwear and rust inhibition properties. In the past, sulfurized triglycerides, such as sulfurized lard oil, have been utilized, especially in association with lightly refined aromatic mineral oils which provided sufficient solubility for the sulfurized triglycerides.\nWith the increased concern for toxicity of aromatic compounds found in such mineral oils, lubricant formulations now comprise essentially non-aromatic oils. This change to substantially non-aromatic base oils created a major problem, resulting from a significant decrease in solubility of the sulfurized triglycerides in the non-aromatic mineral oil, resulting in solidification and/or dropout of the sulfurized triglycerides.\nWhile the solubility problem has been overcome, the modified lubricant products have been found to be either deficient in desirable lubricant properties or incapable of providing needed improvement in these properties.\nIn a typical approach to this problem, as reported in U.S. Pat. No. 3,455,896, sulfurized, low molecular weight polybutenes were reacted with liquid triglycerides, which were susceptible of sulfurization, to yield an additive. In U.S. Pat. No. 3,850,825, another additive was prepared by the sulfurization of a mixture of prime burning lard oil and alkyl oleate. In U.S. Pat. No. 3,740,333, C.sub.10 -C.sub.16 alcohol esters of unsaturated fatty acids, having 18 to 22 carbon atoms, were blended with a triglyceride and either used \"as is\" or sulfurized. Modifications of such compositions have been reported in U.S. Pat. Nos. 4,149,982, 4,166,795, 4,166,796, 4,166,797 and 4,188,300. The use of polymer acids in modified triglycerides for use in soluble oils (aqueous dispersions) has been reported in U.S. Pat. No. 4,067,817.\nAlthough these prior art efforts have increased the solubility of sulfurized fatty oils to acceptable values, there has remained a serious need for sulfurized additives possessing both good solubility and a combination of improved lubricant properties, such as, for example, better low temperature flow properties, better load carrying and antifriction properties, and a lack of sludging. Such improved lubricant properties would also be attractive for use in various fuels systems employed for power generation and heating purposes."}
{"text": "1. Field of the Invention\nThe present invention relates to a fetching method, and more particularly, to a fetching method for accurately fetching the identifications of the SCSI (Small Computer Systems Interface) equipment even under the circumstances interrupted by ESD (Electro-Static Discharge).\n2. Description of the Prior Art\nNowadays, there are many computer equipment such as hard disks and CD-ROMs (Compact Disk Read Only Memory) following the SCSI standard. However, each of the SCSI equipment needs a unique identification for recognizing and fetching their information in a computer system.\nThere usually appears unexpected ESD (Electro-Static Discharge) in the natural world, which usually generates large momentary current. FIG. 2 illustrates the relationship between time and the momentary current caused by the ESD, and the FIG. 3 lists four ESD levels of the momentary currents. It is obvious that the momentary current will rise to the first peak in a range about 0.7 to 1.0 ns (Nano second) whatever its level is, and the momentary current falls down as time passing. For example, an ESD with 4 KV (Kilo Volt) will generate a momentary current with the first peak of about 15 amperes (point A in the FIG. 2). Furthermore, when the time interval starts from the ESD appearance moment to 30 ns (point B in the FIG. 2) and 60 ns (point C in the FIG. 2) later, the momentary current falls down to 8 and 4 amperes, respectively. Please note that the momentary current increases as the voltage of the ESD becomes large.\nThose sensitive computer systems are usually interrupted by the ESD, and thus cause unexpected failures. For example, while the computer system fetches the identification of SCSI equipment, a wrong SCSI identification may be fetched under the situation influenced by the ESD. To overcome this disadvantage, a conventional scheme is to apply a shield to cover the SCSI equipment and therefore redirects the momentary current to a proper place (e.g., to ground). Unfortunately, the shield can not always prevent the SCSI equipment from influencing by the ESD especially when the shield is broken. A requirement has been arisen to disclose a method for accurately fetching the identification of the SCSI equipment anytime."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of soybean breeding. In particular, the invention relates to the novel soybean variety A1024198.\n2. Description of Related Art\nThere are numerous steps in the development of any novel, desirable plant germplasm. Plant breeding begins with the analysis and definition of problems and weaknesses of the current germplasm, the establishment of program goals, and the definition of specific breeding objectives. The next step is selection of germplasm that possess the traits to meet the program goals. The goal is to combine in a single variety an improved combination of desirable traits from the parental germplasm. These important traits may include higher seed yield, resistance to diseases and insects, better stems and roots, tolerance to drought and heat, better agronomic quality, resistance to herbicides, and improvements in compositional traits.\nSoybean, Glycine max (L.), is a valuable field crop. Thus, a continuing goal of plant breeders is to develop stable, high yielding soybean varieties that are agronomically sound. The reasons for this goal are to maximize the amount of grain produced on the land used and to supply food for both animals and humans. To accomplish this goal, the soybean breeder must select and develop soybean plants that have the traits that result in superior varieties."}
{"text": "The present invention is directed to a system for cleaning conveyor belts, and particularly to a cleaning assembly for dual conveyor belt systems.\nA conveyor system employing an endless belt may, over time, require cleaning of the belt surface. This is particularly true when the belt is used for conveying an article during processing of the article where portions of the article, in one form or another, are transformed. Such transformations may be accomplished by, for example, applying pressure or heat to the article being conveyed.\nOne form of such a conveying system employs two endless belts. Each belt traverses a separate belt travel path, with the two belt travel paths together defining an article transport path between opposed faces of the belts. The article being conveyed, upon entry into the article transport path, is simultaneously engaged by the opposed faces of the two belts as it moves along an article processing path. One specific example of such a process is a cooking process, where the article being conveyed is a food item which is cooked while simultaneously being conveyed by the conveyor belt system. In an article cooking application, heat may be applied to the article as it is conveyed by the two belts along the article transport path defined by the belts.\nAs the article is heated, fluids in the article may leave the article in liquid form (i.e., water or oils) or in gas form (i.e., steam or vapor). In addition, portions of the article may separate from the article during processing, causing debris to remain on the belt once the article has left the article transport path. As the endless belts operate over time to transport a plurality of articles, the build up of liquids and debris on the belts can become significant, affecting the performance of the belts in processing the article. In addition, the migration of fluids and/or debris into the operating equipment for the belts and other associated processing equipment can degrade the performance of such equipment. Therefore, it is desirable that any debris or excess liquids created on the belt during processing be removed from the belt in an efficient and expeditious manner."}
{"text": "Voice verification is the process of verifying a person's claimed identity by analyzing a sample of that person's voice. This form of security is based on the premise that each person can be uniquely identified by his or her voice. The degree of security afforded by a verification technique depends on how well the verification algorithm discriminates the voice of an authorized user from all unauthorized users.\nIt would be desirable to use voice verification schemes to verify the identity of a telephone caller. Such schemes, however, have not been successfully implemented. In particular, it has proven difficult to provide cost-effective and accurate voice verification over a telephone network. Generally, this is because the telephone network is a challenging environment that degrades the quality of speech through the introduction of various types of noise and band-limitations. The difficulty in providing telephone-based voice verification is further complicated by the fact that many types of microphones are used in conventional telephone calling stations. These microphones include carbon button handsets, electret handsets and electret speaker phones. Each of these devices possesses unique acoustic properties that affect the way a person's voice may sound over the telephone network.\nGiven the inherent limitations of the prior art as well as the poor frequency response of the telephone network, it has not been possible to successively integrate a voice recognition and verification system into a telephone network."}
{"text": "Sebaceous glands secrete sebum which contain lipids which collect on the skin and in the ear canal. The rate of secretion depends on several factors. Several skin disorders, including acne and seborrhea, are associated with inappropriate sebum production.\nThe removal of excess lipids from the skin can usually be accomplished by bathing using detergents and cleansers. Many preparations have been developed to assist in removal of excess lipids. U.S. Pat. No. 5,026,551 discloses compositions comprising carbon dioxide gas and emulsifiers with oils in bath preparations. In one embodiment the carbon dioxide gas was carried on cyclodextrin. U.S. Pat. No. 4,970,072 discloses use of whey products in bath preparations.\nCyclodextrins have been used as delivery agents for water-insoluble drugs for topical, oral and parenteral delivery. They have also been used to deliver cosmetic preparations to the skin. Several compositions utilize cyclodextrin inclusion products. European patent application 0 366 154 (1990) discloses several cyclodextrin inclusion products for use in cosmetic compositions. U.S. Pat. No. 4,678,598 discloses and claims a cyclodextrin-containing shampoo containing menthol and camphor. The cyclodextrin is provided to depress odor. U.S. Pat. No. 4,267,166 discloses use of cyclodextrin to treat foul breath. U.S. Pat. No. 4,891,361 discloses use of a kojic/cyclodextrin inclusion complexes to prevent elastosis in an animal test by preventing formation of melanin. A publication of Wacker Chemicals discloses that empty cyclodextrins in powder or creams may contain excreted matter of the skin or reaction substances produced on the skin. The statement by Wacker is under a subheading, \"masking of disagreeable smells\" and contains four other statements, all of which pertain to odorous substances. The sebaceous skin lipids are not odorous. Therefore, is appears that the Wacker publication was referring to the containment of unpleasant odors that are produced when bacteria act upon lipids and other skin secretions. There is no evidence therein the materials were useful as cleansing agents for delipidation of the skin."}
{"text": "The present invention is directed to integrated circuits and, more particularly, to a capacitance sensor with noise rejection.\nCapacitance sensors can be used to detect or measure variation of capacitance such as in a capacitive touch panel interface. Common examples of capacitive interfaces are touchpads and touchscreens. In a capacitive touch interface, proximity of an electrically conductive object, such as a human finger or conductive stylus, alters electric fields in the touch interface. The capacitance sensor can sense or measure the physical position or movement of the object relative to the interface and provide a corresponding analog or digital code output signal.\nCapacitive interfaces have a matrix of capacitive elements. Sensing the variations of capacitance in the interface is subject to perturbations, referred to as noise, which may be due to conduction from a power supply or adjacent circuit elements, or induced by external radiation. The noise may be at frequencies higher or lower than the operating frequency of the capacitance sensor.\nIt would be advantageous to have a capacitance sensor with a high level of noise rejection, suitable especially for use in noisy environments."}
{"text": "The present invention relates generally to devices for selectively admitting ventilation air into the interior of a vehicle, and more particularly, to devices of the type which are mounted in a wall of the vehicles and which are operable to selectively open and close the ventilating device.\nMost conventional passenger automobiles include a standard ventilating system that is selectively operable to admit ventilating air that passes through the passenger automobile in a predetermined flow pattern that results in the fresh outside air adequately ventilating the entire interior passenger space within the vehicle. However, in larger vehicles that are designed to carry a relatively large number of passengers, such as a school buses, airport vans, and similar larger passenger vehicles, conventional ventilation systems may not be sufficient to properly ventilate the much larger passenger space within the vehicle. This problem can be exacerbated in special circumstances that further increase the need for more ventilating air, such as larger vehicles that may not include air conditioning and are often required to sit motionless in hot sunlight for prolonged periods of time so that the inside of the vehicle becomes uncomfortably hot.\nTo deal with this problem, it is known that school buses and the like may be provided with ventilating devices which are mounted in the roof of the vehicle, and which can be selectively opened and closed by the operator of the vehicle to admit ventilating air into the vehicle when desired.\nIn one form, such known ventilating devices are operable by an electric motor to tilt upwardly the front edge and/or the rear edge of a ventilating panel, depending upon the desired flow pattern of the ventilating air. If the front edge only is tilted upwardly, fresh air from the outside is forced into the interior of the vehicle by the forward motion of the vehicle. If the rear edge only is tilted upwardly, the forward motion of the vehicle will tend to draw air outwardly from the interior of the vehicle and set-up a flow of ventilating air through the interior of the vehicle. If both the front and rear edges are raised at the same time, the entire panel is moved upwardly to a position generally parallel to it initial position to provide the maximum opening for permitting ventilating air to flow inwardly and/or outwardly through such opening.\nWhile these known ventilating devices generally provide an adequate flow of ventilating air through the vehicle, they have the disadvantage of being somewhat difficult and time consuming to operate. More specifically, the moveable ventilating panel is operated by a single electric motor, and the panel must be moved through a predetermined sequence of steps, normally including four steps, namely:\n(1) tilting the front edge of the panel upwardly;\n(2) then titling the back edge of the panel upwardly;\n(3) then moving the front edge of the panel downwardly to its initial position;\n(4) moving the back edge of the panel downwardly so that the panel is returned to its initial, fully-closed position.\nThus, to get to a particular position of the panel, (e.g. position (3) above, with just the rear edge of the panel tilted upwardly), it may be necessary for the vehicle operator to first move the panel through several positions in the sequence (e.g. positions (1) and (2) above), all of which is time consuming for the operator and which imposes unnecessary wear on the motor and the operating mechanism.\nBy contrast, the present invention provides a ventilating device which overcomes the disadvantages of known ventilating devices and offers significant versatility in the operation of the device.\nIn accordance with the present invention, a selectively openable ventilation device is provided for use in vehicles which includes a frame mountable in a wall of the vehicle and being formed with a ventilation opening therein, and a panel having a configuration corresponding generally to the ventilation opening and having two opposed sides. A first linkage assembly connects one side of the panel to the frame for permitting movement of the one side of the panel between a first closed position at which the one side is disposed within the frame to close the ventilation opening thereat and a second open position at which the one side is disposed in a raised position above the frame to open the ventilation opening thereat. A second linkage assembly connects the other the side of the panel to the frame for permitting movement of the other side of the panel between a first closed position at which the other side is disposed within the frame to close the ventilation opening thereat and a second open position at which the other side is disposed in a raised position above the frame to open the ventilation opening thereat. A first motor is operatively connected to the first linkage assembly to selectively move the one side edge of the panel between the first and second positions thereof, and a second motor is operatively connected to the second linkage assembly to selectively move the other side edge of the panel between the first and second positions thereof, the second motor being operable independently of the first motor whereby either or both of the sides of the panel can be selectively moved toward and way from the frame to permit ventilating air to flow through the ventilation opening at either or both of the sides of the panel.\nIn one preferred embodiment of the present invention, the first and second motors are electric motors, each operated by a control switch mountable within the vehicle and operable to selectively move the two sides of the panel between their open and closed positions and to selectively stop such movement at any desired intermediate position between the open and closed positions.\nAlso, it is preferable that the first and second linkage assemblies each include a pair of spaced mounting brackets having downwardly extending slots therein, with a crossbar positioned in the slots and connected to the motor to be moved in the slots by the motor. A pair of operating links are provided, each operating link being pivotally connected at one of its ends to the crossbar and at the other of its ends to the frame of the ventilation unit, whereby when the crossbar is moved downwardly in the slots by the motor the operating links cause the one panel side to be moved from the first closed position thereof to the second open position, and when the crossbar is moved upwardly in the slots by the motor the operating links cause the one panel side to be moved from the second open position thereof to the first closed position thereof.\nThis preferred operating arrangement may also provide at least one of the operating links with a pin extending outwardly therefrom, and provide a latching member mounted for movement between a first latching position at which it engages the pin to maintain the one panel at its first closed position and a second release position at which it disengages the pin to permit movement of the pin in the slot in the bracket.\nThe latching member is preferably mounted on the bracket for sliding movement relative thereto between the first latching position and the second release position, and it is formed with a hook shaped surface positioned to engage the pin during the last portion of the aforesaid upward movement of the crossbar in the slot. The crossbar may engage the latching member to move the latching member in a direction that will cause the hook shaped surface thereof to urge the pin upwardly when the latching member is in the first latching position and thereby maintain the panel securely pressed against the frame at the first closed position thereof to seal the ventilating opening thereat."}
{"text": "In an LTE wireless communication system, a modulo operation specified by Equation 1 is used to generate a preamble signal, which is transmitted through the physical random access channel (PRACH) for synchronization acquisition.(a·b)mod N (where, 0<a,b≦N)  <Equation 1>\nIn the related art, a linear search scheme is employed to physically implement the modulo operation in Equation 1. In this scheme, using comparison between the dividend at and the divisor N, subtraction and comparison operations are repeated until the value of (a·b)·N becomes less than the value of N, and the final result becomes the remainder. However, utilization of such a linear search scheme for the modulo operation of Equation 1 necessitates employment of a divider, which may increase hardware complexity and worsen synthesis timing.\nIn performing the modulo operation using the above linear search scheme, as the amount of computation is linearly proportional to the magnitude of the dividend a·b, many stages of computation are needed to obtain the remainder and hence computation time becomes longer. In addition, as the divisor N becomes larger, the above related art method becomes less effective."}
{"text": "An embedded universal integrated circuit card (eUICC) is a traditional universal integrated circuit card (UICC) embedded in a terminal, but is not removable, and therefore cannot be arbitrarily replaced like a traditional UICC card. However, a profile in the eUICC card can be changed through remote management. The profile is a combination of a file structure, an application, data, and the like, and after the profile is configured in the eUICC, the eUICC may access a mobile network by using the profile. Specifically, a profile manager (PM) outside the eUICC card may manage all profile management functions of each profile in the eUICC card by using a corresponding management function in the eUICC card, such as a profile management domain (PMD). The profile management functions may include: transmission, activation, deactivation, deletion, and the like. However, if the PM is used to manage the profile in the eUICC card, the PM needs to be authorized.\nIn some approaches, authorization information is stored in an authorization server, and each time after a communications device is started, the communications device needs to report a device serial number to the authorization server, to obtain the authorization information from the authorization server. However, efficiency of the authorization manner is low."}
{"text": "1. Field of the Present Disclosure\nThe present disclosure relates to an image forming apparatus.\n2. Description of the Related Art\nIn an image forming apparatus such as a printer, a copy machine, a facsimile machine or a multi function peripheral, calibration is performed as a density correction technique for image forming. In the calibration, light is irradiated to a photo conductor drum by an exposure device and consequently an electrostatic latent image of test patches (an image for density correction) is formed on the photo conductor drum; a toner image is formed by supplying toner to the electrostatic latent image from a development device; and after transferring the toner image to an intermediate transfer drum, densities of the image is optically measured using an optical sensor or the like, exposure light intensity of the exposure device, a surface voltage of the photo conductor drum, a development bias of the development device and so forth are adjusted on the basis of the measured density of the image for correcting densities of image forming.\nIn an image forming apparatus that includes a development roller arranged in a development device in order to supply toner to a photo conductor drum, if development characteristic is not uniform to rotation angles of the development roller due to eccentricity of a rotation shaft, unevenness of a roller surface of the development roller and so forth, then developer is ununiformly distributed on the outer surface of the development roller, toner does not precisely adhere from the development roller to a photo conductor drum by an expected toner amount, and consequently a whole area of a test patch may not have an uniform density on the photo conductor drum. In such a case, the precise calibration may not be performed.\nIn an image forming apparatus, in order to reduce density unevenness of a test patch, a test patch is used of which a length on a photo conductor drum is larger than a periphery length of a development roller, and an average value of densities within the test patch is measured.\nIn the aforementioned image forming apparatus, the length of the test patch is larger than the periphery length of the development roller, and therefore if the periphery length of the development roller is large then the length of the test patch must be large. A large length of the test patch results in a large toner consumption amount on the test patch and a long time required for the calibration.\nFurther, for example, it may be considered that two pattern images that are two sets of test patches corresponding to densities in gradation are formed and a measured density corresponding to a density setting value is calculated from a measured density of a test patch of the first set and a measured density of a test patch of the second set. However, as shown in FIG. 6A and FIG. 6B, if a position of the test patch of the first set and a position of the test patch of the second set within a fluctuation pitch (i.e. a fluctuation period) of density unevenness are close to each other, the density unevenness still affects the measured densities of the both test patches and consequently influence of the density unevenness is not reduced."}
{"text": "1. Field of the Invention\nThis invention relates generally to acid baths, and more particularly to high temperature acid baths as used for etching semiconductors and the like.\n2. Description of the Prior Art\nDuring the manufacture of integrated circuits and like semiconductor devices, it is generally necessary to place in a bath of a hot acid a usually rectangular rack known as a boat, holding one or more disks, called slices or wafers, of semiconductor material, usually pure silicon, on which a large number of integrated circuits have been formed by a known process using a material referred to as a \"photoresist\" in order to etch, or wash, away unwanted material on the disks. The disks subsequently are scored and broken into individual integrated circuit chips.\nThe container of the hot acid, which container is generally referred to as a beaker, is disposed over a sink or other drain, so to as much as possible prevent the acid, which is extremely toxic and highly volatile, from escaping the area of the bath. For this reason, it is imperative that a good seal be obtained between a peripheral portion of the open-top acid beaker and a structure supporting the beaker. Otherwise, acid and its fumes will collect beneath the beaker resulting in corrosive and possibly explosive conditions.\nAlthough for safety reasons the acid bath usually is placed over some sort of drain, as opposed to what might be referred to as \"free standing\", this alone generally is not sufficient to remove the spilled acid and vapors rapidly enough. A principal difficulty arises in creating the aforementioned seal between an acid beaker or vessel and its support structure due to the fact that the vessel, which is generally constructed from either quartz or a borosilicate glass, is an open-topped device of rectangular configuration. While the rim area of round beakers constructed from quartz or glass can easily be molded into a configuration permitting a good seal with an associated support structure, it has been found to be very difficult to achieve a suitable seal with a rectangular beaker configuration. Round beaker, however, are not desirable for use in the semiconductor processing industries because such beakers waste much of the volume of the bath since the rectangular racks, or boats, will fit in only a small portion of the horizontal section of the beaker.\nOne prior art approach to achieving such a seal provided a horizontal flange on the upper rim of a vessel so as to be supported by a shelf forming part of the associated support structure. Inserted between the flange and the shelf was a room-temperature vulcanizing (RTV) elastomer for effecting the desired seal between the elements. A principal drawback encountered with this construction was that it permitted acid to collect immediately adjacent the seal, with the result that the acid eventually worked its way through the elastomer and permitted leakage by the seal, creating corrosive and possibly explosive conditions."}
{"text": "Well completions are used in a variety of well related applications involving, for example, the production or injection of fluids. Generally, a wellbore is drilled, and completion equipment is lowered into the wellbore by tubing or other deployment mechanisms. The wellbore may be drilled through one or more formations containing desirable fluids, such as hydrocarbon based fluids.\nIn applications in which the wellbore has been formed through a plurality of formations, the wellbore often is divided into wellbore zones to better control the flow of fluid between each formation and the wellbore. Accordingly, it can be beneficial to have at least some control over the production of fluid from individual formations and/or over the injection of fluid into individual formations. The completion equipment may comprise devices, such as packers and multiple pumps, that can help control fluid flow with respect to each formation. However, the ability to efficiently control fluid flow in such subterranean environments while monitoring well conditions can be difficult."}
{"text": "Attaching accessories to a helmet can yield the helmet unsafe to wear. For example, a camera attached to a helmet may penetrate through the helmet and cause significant damage to the user. The helmet can become also dangerous if the accessory is mounted rigidly so that in an event of an impact to the accessory, the impact will cause tangential force or torque to user's head and can thus cause severe damages.\nSome attempted solutions exist which utilize mounting platforms, base or rails attached to the surface of the helmet on which the accessory such as camera has been attached. These solutions can make the use of camera more comfortable but doesn't prevent the accessory to cause danger.\nSome attempts utilize magnets for attaching the accessories to the helmet. While the accessory can detach from the magnet quite easily in case of an impact, the magnet itself is still rigidly attached to the helmet and can become dangerous.\nThere is still need for a helmet to which accessories can be attached without jeopardizing the safety of the helmet."}
{"text": "The present invention relates to a lighting device and associated circuitry for lighting and dimming a semiconductor light-emitting element such as a light emitting diode (LED) or a plurality of the same forming an LED module. More particularly, the present invention relates to an LED lighting device effective to provide a dimming operation based on externally provided PWM dimming signals, and illumination fixtures or video cameras utilizing the same.\nAn exemplary LED lighting device as conventionally known in the art, and as represented in FIG. 7, may include a switching element Q1 which is coupled in series to a direct-current (DC) power source 2, and an inductor L1 coupled in series with the switching element Q1 and to which an current flows from the DC power source 2 when the switching element Q1 is turned on. Energy stored in the inductor L1 (or equivalent inductive element) when the switching element Q1 is turned on is discharged to a semiconductor light-emitting element 4 via regeneration diode D1 when the switching element is turned off. A current detector R1 detects an current flowing to the switching element Q1. A control circuit is configured to turn the switching element Q1 off when the detected current reaches a predetermined value (i.e., an on-voltage of transistor Tr1) and to turn the switching element Q1 on when the energy discharge of the inductor L1 is completed (when a diode D2 is turned off). However, a LED lighting device so configured does not typically support a dimming function.\nFIG. 13 shows an example of a conventional LED lighting device which is capable of supporting a dimming function and which (a) controls a switching element Q1 to be turned off when an current flowing to an inductor L1 reaches a predetermined value, and (b) controls the switching element Q1 to be turned on when an current is fully discharged from the inductor L1 to a semiconductor light-emitting element 4 via a regeneration diode D1. The lighting device controls the average value of the current flowing to the LED 4 to be constant based on a reference voltage Vref1 by improving the input power factor. Dimming of the LED 4 can also or alternatively be controlled by adjusting the reference voltage Vref1. A lighting device so configured can control the average current flowing to the semiconductor light emitting element 4 to be a constant value based on the reference voltage Vref1 even when a power source voltage and an ambient temperature have changed, and can accordingly reduce input current distortion.\nSince the control circuit can be configured by using a commercially available IC for power factor improvement, it can be produced at relatively low cost. Among the commercially available ICs for power factor improvement are low-cost ICs which include an error amplifier EA, a multiplier circuit 52, a comparator CP1, a flip-flop FF1, and a driving circuit 54 in one chip. However, as previously noted, the semiconductor light-emitting element 4 is dimmed by setting the reference voltage Vref1 to be adjustable, and dimming control is substantially more difficult where the reference voltage Vref1 is incorporated in the IC.\nAnother conventional dimming method includes converting an externally-provided PWM signal into a second PWM signal having a different pulse width, interrupting current flowing to a semiconductor light-emitting element in accordance with the converted PWM signal, and adjusting the amplitude of the current flowing to the semiconductor light-emitting element in accordance with a direct-current voltage obtained by smoothing the PWM signal. In such techniques, the on-resistance of the transistor is variably controlled to adjust the amplitude of the current flowing to the semiconductor light-emitting element, and thus a large loss of electric power is generally incurred. To reduce power loss, the amplitude of the current flowing to the semiconductor light-emitting element is adjusted by using a switching power source circuit such as a chopper circuit, which has high power conversion efficiency when operating in a critical mode where the switching element is controlled to be turned on after detecting a zero-cross point for the regeneration current.\nAn exemplary configuration for controlling the current flowing to the semiconductor light-emitting element to be constant with use of a step-down chopper circuit operating in the critical mode includes PWM-controlling a drive signal of a switching element of the step-down chopper circuit on the basis of a dimming signal.\nHowever, when the duty cycle (on-pulse width) of the switching element is controlled by the PWM, the zero-crossing time for the regeneration current, and thus the switching frequency, may vary in a wide range when control of the duty cycle is arbitrarily variable on the basis of the dimming signal in the switching power source circuit operating in the critical mode, as shown for example by the dashed line in FIG. 15(e)). Also, even if the high frequency switching operation is intermittently stopped on the basis of a low-frequency PWM signal, dimming cannot be sufficiently realized only within the range where the PWM signal can be varied, and accordingly the dimming range is limited.\nReferring now to FIG. 19, a conventional operating method for an LED lighting device is represented wherein current flowing to a semiconductor light-emitting element (or LED module) is controlled via a step-down chopper circuit operating in a so-called critical mode. Energy stored in an inductor during an on-period TON of the switching element is discharged during an off-period TOFF, the switch is turned on again when the energy discharge is completed, and thus the power conversion efficiency becomes greater in comparison with other control modes. In addition, since substantially half the peak value of the switching current is an effective value of a load current, constant current control can be easily realized.\nAn example of the critical mode operation may be further described with reference to a switching element Q1 of a step-down chopper circuit 3a as shown in FIG. 5(a). A DC voltage is supplied across input terminals A and B, having conventionally been obtained by stepping-up a commercially-available alternating-current (AC) power source. An LED series circuit or a load circuit is coupled across output terminals C and D. When the switching element Q1 is turned on, current IQ1 (referring again to FIG. 19) flows via a switching element Q1, an inductor L1, and a capacitor C2, and thus energy is stored in the inductor L1. When the switching element Q1 is turned off, a counter-electromotive force is generated due to the energy stored in the inductor L1, and thus a regeneration current ID1 flows from the inductor L1 to the capacitor C2 and then to a diode D1. When the switching element Q1 is turned on again (upon the regeneration current ID1 returning to zero), switching loss is relatively small and accordingly the power conversion efficiency becomes higher in comparison with other control modes.\nIt would be desirable to provide a circuit configuration for PWM control of the on-period TON of the switching element Q1 based on an external dimming signal, and which could be implemented prior to factory shipment via easy and low-cost output adjustments to variations in LED characteristic variations or circuit constants such as inductance, temperature variations, age deterioration, etc."}
{"text": "Augmented reality allows interaction among users, real-world objects, and virtual or computer-generated objects and information within an environment. The environment may be, for example, a room equipped with computerized projection and imaging systems that enable presentation of images on various objects within the room and facilitate user interaction with the images and/or objects. The augmented reality may range in sophistication from partial augmentation, such as projecting a single image onto a surface and monitoring user interaction with the image, to full augmentation where an entire room is transformed into another reality for the user's senses. The user can interact with the environment in many ways, including through motion, gestures, voice, and so forth.\nTo enable such augmented reality environments, however, there is a continuing need for improved projection systems. Such improvements might include lighter weight, smaller form factors, and less obtrusive integration into the environment."}
{"text": "Japanese Patent Laying-open No. 2003-304604 discloses a device driving a motor mounted as a source driving electric vehicles, hybrid vehicles, fuel cell vehicles and other various types of vehicles.\nIn this device, if a temperature detection means detects that the motor has a temperature equal to or higher than a temperature limit, a motor output control means limits the output of the motor. Herein a means for detecting a rate of change in temperature detects the rate of change in temperature of the motor and modifies a setting of the temperature limit in accordance with the rate of change detected.\nMore specifically, if the motor's rate of change in temperature is equal to or larger than a predetermined rate of change, the means for detecting a rate of change in temperature determines that the motor has a large increase in temperature, and accordingly the means sets the current temperature limit to be a first temperature limit, and when the motor attains the first temperature limit or higher the motor output control means limits the output of the motor.\nIn contrast, if the motor's rate of change in temperature is smaller than the predetermined rate of change, the means for detecting a rate of change in temperature determines that the motor has a small increase in temperature, and accordingly the means sets the current temperature limit to be a second temperature limit higher than the first temperature limit, and when the motor attains the second temperature limit or higher the motor output control means limits the output of the motor.\nThis device allows a vehicle to run an increased distance without limiting the output of the motor, can protect the motor from high temperature, and also allows the motor to fully exhibit its performance.\nHowever, the motor drive device disclosed in Japanese Patent Laying-open No. 2003-304604 only limits the output of the motor in accordance with the rate of change in the temperature of the motor in view of protecting the motor from high temperature. It does not give consideration to the state of the vehicle.\nFor example, when a vehicle is in a state in which while it is running up hill it may fall downhill, such state should first and most of all be avoided. The above described motor drive device, despite the vehicle's such state, may prioritize protecting the motor from high temperature and stringently limit the motor's output. Thus with the above described motor drive device there is a possibility that while the vehicle is running uphill it may fall downhill."}
{"text": "1. Technical Field\nThe present invention relates generally to semiconductor technology and more specifically to semiconductor research and development.\n2. Background Art\nToday, electronic products are used in almost every aspect of life, and the heart of these electronic products is the integrated circuit. Integrated circuits are used in everything from airplanes and televisions to wristwatches.\nIntegrated circuits are made in and on silicon wafers by extremely complex systems that require the coordination of hundreds or even thousands of precisely controlled processes to produce a finished semiconductor wafer. Each finished semiconductor wafer has hundreds to tens of thousands of integrated circuits, each worth hundreds or thousands of dollars.\nThe ideal would be to have every one of the integrated circuits on a wafer functional and within specifications, but because of the sheer numbers of processes and minute variations in the processes, this rarely occurs. “Yield” is the measure of how many “good” integrated circuits there are on a wafer divided by the total number of integrated circuits formed on the wafer divided by the maximum number of possible good integrated circuits on the wafer. A 100% yield is extremely difficult to obtain because minor variations, due to such factors as timing, temperature, and materials, substantially affect a process. Further, one process often affects a number of other processes, often in unpredictable ways.\nIn a manufacturing environment, the primary purpose of experimentation is to increase the yield. Experiments are performed in-line and at the end of the production line with both production wafers and experimental wafers. However, yield enhancement methodologies in the manufacturing environment produce an abundance of very detailed data for a large number of wafers on processes subject only to minor variations. Major variations in the processes are not possible because of the time and cost of using production equipment and production wafers. Setup times for equipment and processing time can range from weeks to months, and processed wafers can each contain hundreds of thousands of dollars worth of integrated circuits.\nThe learning cycle for the improvement of systems and processes requires coming up with an idea, formulating a test(s) of the idea, testing the idea to obtain data, studying the data to determine the correctness of the idea, and developing new ideas based on the correctness of the first idea. The faster the correctness of ideas can be determined, the faster new ideas can be developed. Unfortunately, the manufacturing environment provides a slow learning cycle because of manufacturing time and cost.\nRecently, the great increase in the complexity of integrated circuit manufacturing processes and the decrease in time between new product conception and market introduction have both created the need for speeding up the learning cycle.\nThis has been accomplished in part by the unique development of the integrated circuit research and development environment. In this environment, the learning cycle has been greatly speeded up and innovative techniques have been developed that have been extrapolated to high volume manufacturing facilities.\nTo speed up the learning cycle, processes are speeded up and major variations are made to many processes, but only a few wafers are processed to reduce cost. The research and development environment has resulted in the generation of tremendous amounts of data and analysis for all the different processes and variations. This, in turn, has required a large number of engineers to do the analysis. With more data, the answer always has been to hire more engineers.\nHowever, this is not a solution for major problems.\nThe problems include, but are not limited to, the requirement that all nonphotolithographic processes currently involve treatment of an entire wafer at one time. Therefore, it is not possible to perform different experiments of nonphotolithographic processes on different areas of the same wafer.\nThe problems include, but are not limited to, evaluating lot-to-lot differences in parametric tests involving generation of multiple sets of data or families of curves. Often, these curves exhibit certain characteristics, such as inflection points, nodes, etc., which are indicative of product performance of certain known failure mechanisms. Unfortunately, evaluation of the curves requires “expert” eyes because the differences can be extremely subtle.\nThe problems include, but are not limited to, the need to statistically collect large amounts of current versus voltage data.\nThe problems include, but are not limited to, improving photolithography focus exposure matrix design comparisons.\nSolutions to these problems have been long sought but prior developments have not taught or suggested any solutions and, thus, solutions to these problems have long eluded those skilled in the art."}
{"text": "Online advertisement service providers (“ad providers”) provide Internet advertisements (“ads”) to website pages. An advertiser places the ad with the ad provider.\nA content network includes a plurality of website pages provided by publishers that provide content (e.g., text, audio, videos, articles, blogs, etc.) (“content pages”) and allow advertisers to show ads along with the content of the content pages. A publisher publishes the content pages (e.g., by making the content pages available to users). Publishers may publish or produce news web sites, weather web sites, sports web sites, game web sites, guitar websites, and/or websites relating to other subjects. Additionally, the web sites can be informational, interactive, commercial, or other types of web sites. Content pages differ from search result pages, which are pages that solely serve to provide search responses to a user query."}
{"text": "Cyclodextrins make up a family of cyclic oligo- and polysaccharides containing 5 or more D-glucopyranoside units linked through a 1-4 glycosidic bonds. The most typical cyclodextrins contain 6 to 8 glucopyranoside units in a ring, creating a cone shape. Within this family, β-cyclodextrin have 7 glucopyranoside units in a ring. Classic sulfation procedures result in the sulfation of one, two or three hydroxyl groups of the glucopyranoside units.\nSome medicinal uses of polysulphated cyclodextrins are already known in the art. For instance U.S. Pat. No. 6,930,098 teaches treating a human afflicted with arthrosis by administering a composition consisting essentially of a polysulphated cyclodextrin or an addition salt thereof, in combination with one or more non-toxic pharmaceutically acceptable excipients; preferably the dosage for such treatment ranges from 50 mg to 1,500 mg per day. According to Osteoarthritis and Cartilage (2008) 16:986-993, β-cyclodextrin polysulfate (CDPS), subcutaneously administered in a rabbit model of experimental osteoarthritis (OA) reduced the cartilage lesions and osteocyte formation in the affected joints. These data suggest that CDPS positively affects the tissue pathology underlying OA and this agent can therefore be classified as a structure or disease modifying OA drug.\nApart from the chondroprotective capacities, other characteristics of CDPS are described in the literature. Polysulfated cyclodextrins have been shown to possess important biological activities similar to those of heparin which can be explained by the similarities in their molecular structure. Apart from a strong binding affinity for a fibroblast growth factor, an anti-angiogeneic activity and a capacity of inhibiting cellular invasion by HIV retro virus, both polysaccharides possess anticoagulant properties and CDPS may elicit heparin-induced thrombocytopenia.\nIn vitro studies concerning CDPS and coagulation show prolongation of thrombin clotting times and a reduction in thrombus formation. Such biological activities were clearly shown to be related to the molecular structure of the polysaccharide and varied following distinct modifications. The introduction of distinct alkyl groups for example, did not prolong the activated partial thromboplastin time (aPTT) in vitro suggesting a reduced anti-coagulant activity compared to CDPS.\nNext, the potency to induce heparin-induced thrombocytopenia (hereinafter referred as HIT) and thromboembolic accidents through cross reaction with heparin/platelet factor IV antibodies is a matter of concern. These antibodies can arise occasionally when activated thrombocytes release platelet factor IV (PF4) during heparin treatment. Heparin then forms a complex with PF4 that acts as an antigen which triggers the production of auto-antibodies. These anti-bodies bind to the complex via their F(ab)′ region and to the FcγRII (IgG CD32) of other platelets via the Fc portion thereby initiating platelet activation, aggregation and generation of platelet-derived microparticles. These pro-coagulant particles are likely to induce the thrombotic complications of HIT. It has been shown that some low molecular weight heparins as well as other sulphated polysaccharides, e. g. chondroitin polysulphates, can also bind to HIT antibodies in the presence of PF4 and that the reactivity is dependant on their molecular weight and the sulfation grade. Only a subset of patients will produce antibodies with platelet-activating properties after heparin therapy. Some of the later will develop thrombocytopenia and even less will develop thrombosis.\nThus, it is one problem to be addressed by the presently claimed invention to provide novel therapeutic agents with a preserved chondro-protective capacity, a reduced effect on coagulation and a reduced risk for heparin-induced thrombocytopenia.\nIt is another problem to be addressed by the presently claimed invention to provide novel therapeutic agents that are suitable for the prevention of platelet aggregation and vascular thrombosis in individuals afflicted with HIT syndrome.\nIt is another problem to be addressed by the present invention to provide novel therapeutic agents which can be used to treat HIT syndrome without inducing anti-coagulant activity.\nFurthermore, it is another problem to be addressed by the presently claimed invention to provide novel therapeutic agents which can be used to treat HIT syndrome without inducing platelet activation or thrombosis in the presence of heparin- and platelet factor IV-complex reactive antibodies.\nIt is another problem to be addressed by the presently claimed invention to provide novel therapeutic agents which can be used for the treatment and/or prophylaxis of degenerative joint diseases such as osteoarthritis, articular rheumatism, arthrosis or degenerative arthritis, or for cartilage repair or connective tissue repair."}
{"text": "Biaxial films having controlled refractive indexes along the directions of three axes are useful in the optical field utilizing polarized light. In particular, in the field of liquid crystal display, such films permitting precise control of polarized light are of a great importance.\nIn the case of preparing such optically biaxial film, it is a popular process to obtain the film by biaxially stretching a film obtained from a polymer (see, for example, JP-A-2-264905). In the case of obtaining a biaxial film by the biaxially stretching technique, the refractive indexes along the directions of three axes can be controlled to desired levels with comparative ease because the refractive indexes can be controlled through draw ratio.\nOn the other hand, it has recently been proposed to obtain a biaxial film by using a biaxial liquid crystal (see, for example, JP-A-2002-6138 and JP-A-2002-174730). In comparison with biaxially stretched films which have conventionally been often used, the biaxial film using the biaxial liquid crystal has a merit of reducing its thickness so much that using the biaxial liquid crystal for the biaxial film is a useful means for reducing the thickness or weight of a device.\nHowever, in the case of preparing a biaxial film by using a biaxial liquid crystal, there arises a problem that refractive indexes of the film along the directions of three axes can not arbitrarily be controlled. This is attributed to that the refractive indexes of a resultant biaxial film along the directions of three axes are almost decisively determined by the refractive indexes of the liquid crystal phase-developing compound (liquid crystal compound) along the directions of three axes. That is, in order to adjust the refractive indexes of the biaxial film along the directions of three axes to desired levels, it is the only method to synthesize a liquid crystal compound having desired refractive indexes. However, in comparison with the number of compounds developing uniaxial liquid crystal, the number of compounds capable of developing biaxial liquid crystal is so small (see, for example, D. Demus, J. Goodby, et al., Handbook of Liquid Crystals Vol. 2B; Low Molecular Weight Liquid Crystals II, WILEY-VCH p. 933-943) that it has been extremely difficult to arbitrarily control the refractive indexes along the directions of three axes."}
{"text": "1. Technical Field\nThe subject matter described here generally relates to static structures including a rigid member having a limited closed periphery and which is greatly elongated relative to any lateral dimension, and, more particularly, to joints for wall sections of wind turbine towers.\n2. Related Art\nA wind turbine is a machine for converting the kinetic energy in wind into mechanical energy. If that mechanical energy is used directly by machinery, such as to pump water or to grind wheat, then the wind turbine may be referred to as a windmill. Similarly, if the mechanical energy is further transformed into electrical energy, then the turbine may be referred to as a wind generator or wind power plant.\nWind turbines use one or more airfoils in the form of a “blade” to generate lift and capture momentum from moving air that is them imparted to a rotor. Each blade is typically secured at its “root” end, and then “spans” radially “outboard” to a free, “tip” end. The front, or “leading edge,” of the blade connects the forward-most points of the blade that first contact the air. The rear, or “trailing edge,” of the blade is where airflow that has been separated by the leading edge rejoins after passing over the suction and pressure surfaces of the blade. A “chord line” connects the leading and trailing edges of the blade in the direction of the typical airflow across the blade. The length of the chord line is simply the “chord.”\nWind turbines are typically categorized according to the vertical or horizontal axis about which the blades rotate. One so-called horizontal-axis wind generator is schematically illustrated in FIG. 1 and available from GE Energy of Atlanta, Ga. USA. This particular configuration for a wind turbine 2 includes a tower 4 supporting a drive train 6 with a rotor 8 that is covered by a protective enclosure referred to as a “nacelle.” The blades 10 are arranged at one end of the rotor 8, outside the nacelle, for driving a gearbox 12 connected to an electrical generator 14 at the other end of the drive train 6 arranged inside the nacelle along with a control system 16. An anemometer 18 is also provided on the nacelle for providing information to the control system 16.\nThe wind turbine tower 4 is typically formed from several tubular, frustro-conical sections that are stacked on top of each other. For example, commonly-owned co-pending U.S. patent application Ser. No. 11/399,829 by Ronald Cairo discloses “Methods and Apparatus for Assembling Wind Turbine Towers.” In that publication, the members include a protrusion that extends radially and circumferentially outward at a bottom end, defining an annular notch for receiving the top end of a second tower member. At least one bolt may be extended through second tower member protrusion and first tower member top end in order to increase the structural integrity of the joint formed by interlocking first tower member and the second tower member. The tower members may also be provided with vertical flanges that are substantially aligned to form a joint.\nOther wind turbine tower configurations have been provided with inner or outer flanges that that are normally welded to the peripheral wall portions. For example, commonly-owned and co-pending U.S. patent application Ser. No. 10/517,585 by Wolfgang Fuelhaas and Holger Lühn discloses a “Method for Generating a Substantially Uninterrupted Connection of the Peripheral Wall Portions,” and an English-language abstract of German Patent Publication No. DE10152018 by Marc Seidel discloses “Component arrangement for manufacturing tower for wind power system has components with boundary surfaces bounding test channel enabling testing of contact between components.”\nHowever, wind turbine towers with such flanged, connections suffer from a variety of drawbacks. For example, the flanges require additional material, welding, inspection and labor resources to install and inspect. Flanges can also cause positioning and bending problems, particularly during transportation and storage where protruding flanges can interfere with stacking and impinge upon other structures."}
{"text": "The present invention relates to management of a multi-tier storage environment, and more specifically, this invention relates to efficient management of high performance tiers in a multi-tier storage environment.\nA file system defines how files are named and manages how they are placed for storage and retrieval. File system functionality may be divided into two components: a user component and a storage component. The user component is responsible for managing files within directories, file path traversals, and user access to files. The storage component of the file system determines how files are stored physically on the storage device.\nIn addition, a file system may attempt to efficiently place data in different locations according to the importance of the data and how frequently it is accessed. When data is migrated from one physical location to another, the strategic placement of the data may be lost and applications which access the data may have their performance suffer."}
{"text": "The normal human heart is a strong, muscular pump a little larger than a fist. It pumps blood continuously through the circulatory system. Each day the average heart “beats” (or expands and contracts) 100,000 times and pumps about 2,000 gallons of blood. In a 70-year lifetime, an average human heart beats more than 2.5 billion times.\nThe heart pumps blood through a circulatory system, which is a network of elastic tubes through which blood flows as it carries oxygen and nutrients to all parts of the body. The circulatory system includes the heart, lungs, arteries, arterioles (small arteries), and capillaries (minute blood vessels). It also includes venules (small veins) and veins, the blood vessels through which blood flows as it returns to the heart.\nThe circulating blood brings oxygen and nutrients to all the organs and tissues of the body, including the heart itself. The blood also picks up waste products from the body's cells. These waste products are removed as they're filtered through the kidneys, liver and lungs.\nOver time, the coronary arteries which supply the heart muscle with blood can become clogged. One cause of clogged arteries is due to a condition called atherosclerosis, or hardening of the arteries. Atherosclerosis causes a constriction of the inner lumen of the affected artery when the lumen of the arteries become more narrow due to a pathological accumulation of cells, fats and cholesterol called plaque. The descriptive term given to this narrowing of the coronary arteries is “stenosis.” Stenosis means constriction or narrowing. A coronary artery that is constricted or narrowed is referred to as stenosed. When stenosis of the coronary artery is sufficient to deprive the heart muscle of the oxygen levels necessary for cell viability, the result is typically myocardial infarction, typically referred to as a heart attack.\nA heart attack occurs when the blood supply to part of the heart muscle itself, the myocardium, ceases or is severely reduced. This occurs when one or more of the arteries supplying blood to the heart muscle (coronary arteries) becomes partially or completely obstructed by plaque stenoses. If cessation of the blood supply occurs for a long time, heart muscle cells suffer irreversible injury and die. Severe disability or death can result, depending on how much heart muscle is damaged.\nCoronary artery bypass surgery is a heart operation used to treat coronary artery disease. In coronary artery bypass surgery a blood vessel is used to go around or “bypass” clogged coronary (heart) arteries. During the “bypass” procedure, a blood vessel from the patient's chest or leg is used as the “bypass” conduit. For venous “bypass” grafts, one end of the vessel is attached to the aorta (the large artery coming out of the heart) and the other end is attached to the coronary artery below the point where it's clogged. Once the clog has been bypassed, blood can once again flow through the bypass graft to the heart, in a manner that prevents ischemia and infarction. Almost half a million coronary bypass operations are performed each year in the USA.\nAnother procedure for opening clogged coronary arteries is to perform percutaneous transluminal coronary angioplasty, or balloon angioplasty. Balloon angioplasty is an established and effective therapy for some patients with coronary artery disease. Balloon angioplasty is used to dilate (widen) arteries narrowed by plaque. During the procedure, a catheter with a deflated balloon on its tip is passed into the narrowed part of the artery. The balloon in then inflated, and the narrowed area is widened. Balloon angioplasty is a less traumatic and less expensive alternative to bypass surgery for some patients with coronary artery disease. However, in 25 to 30 percent of patients the dilated segment of the artery renarrows (restenosis) within six months after the procedure. The patient may then require either to repeat the balloon angioplasty or to undergo coronary bypass surgery.\nOne approach to preventing restenosis has been to insert a “stent” across the stenosed area of coronary artery. A stent is a metallic wire mesh tube that is used to prop open an artery that has been recently dilated using balloon angioplasty. The stent is collapsed to a small diameter, placed over an angioplasty balloon catheter and moved into the area of the blockage. When the balloon is inflated, the stent expands, locking in place to form a rigid support (structural scaffolding) which holds the artery lumen open. The stent remains in the artery permanently to help improve blood flow to the heart muscle. However, reclosure (restenosis) remains an important issue with the stent procedure.\nSeveral approaches have been taken to reduce the occurrence of restenosis associated with the stent procedure. Stents have been impregnated with drugs and chemicals that emit radiation (gamma-rays) in an attempt to reduce the frequency of restenosis. Also, drug eluting stents have been used in an attempt to reduce the occurrence of restenosis. However, a need still exists for additional safe and effective treatments to prevent restenosis after the placement of an intravascular stent."}
{"text": "Phytate is the major storage form of phosphorus in cereals and legumes. However, monogastric animals such as pig, poultry and fish are not able to metabolise or absorb phytate (or phytic acid) and therefore it is excreted, leading to phosphorous pollution in areas of intense livestock production. Moreover, phytic acid also acts as an antinutritional agent in monogastric animals by chelating metal agents such as calcium, copper and zinc.\nIn order to provide sufficient phosphates for growth and health of these animals, inorganic phosphate is added to their diets. Such addition can be costly and further increases pollution problems.\nThrough the action of phytase, phytate is generally hydrolysed to give lower inositol-phosphates and inorganic phosphate. Phytases are useful as additives to animal feeds where they improve the availability of organic phosphorus to the animal and decrease phosphate pollution of the environment (Wodzinski R J, Ullah A H. Adv Appl Microbiol. 42, 263-302 (1996)).\nA number of phytases of fungal (Wyss M. et al., Appl. Environ. Microbiol. 65 (2), 367-373 (1999); Berka R. M. et al., Appl. Environ. Microbiol. 64 (11), 4423-4427 (1998); Lassen S. et al., Appl. Environ. Microbiol. 67 (10), 4701-4707 (2001)) and bacterial (Greiner R. et al Arch. Biochem. Biophys. 303 (1), 107-113 (1993); Kerovuo et al., Appl. Environ. Microbiol. 64 (6), 2079-2085 (1998); Kim H. W. et al., Biotechnol. Lett. 25, 1231-1234 (2003); Greiner R. et al., Arch. Biochem. Biophys. 341 (2), 201-206 (1997); Yoon S. J. et al., Enzyme and microbial technol. 18, 449-454 (1996); Zinin N. V. et al., FEMS Microbiol. Lett. 236, 283-290 (2004))) origin have been described in the literature."}
{"text": "Sleep disorders are extremely common in the general population, with estimates of 15 million adults in the U.S. with persistent insomnia, and another 24 million adults and children suffering from obstructive sleep apnea. The current model of care is for people to seek consultation with their primary care doctor first about their sleep issues. From there, people are referred for sleep studies, sleep specialists, or given medications. Treatment is often delayed, involves potentially habit-forming if medications, or is too expensive to pursue.\nSmartphones are becoming the predominant link between data and people. Most current smartphones provide a capability to use mobile software applications (apps). A mobile software application (app) can embody a defined set of functionality and can be installed and executed on a mobile device, such as a smartphone, a tablet device, laptop computer, or other form of mobile computing or communications device. Conventional mobile apps are available that focus on one narrow aspect or another of different sleep problems. However, these conventional mobile apps do not integrate a set of diagnostic and therapeutic tools to fully assess and treat a user's sleep disorder. For example, current systems do not integrate a user's sleep history in a diagnostic and treatment system. Additionally, one may find a mobile app that plays soothing music, or another that wakes the user up at a particular time, or another that records snoring. However, these systems cannot combine data obtained in a diagnostic phase for use in a sleep disorder treatment phase."}
{"text": "In recent years, there has been enormous increase in the amount of digital content available online. As a result, the amount of content downloaded by users has increased. This content includes audio and video content. For example, Apple® iTunes® enables users to legally purchase and download music to personal music devices. Video services exist that allow users to access and download video media. Due to the large amount of content available, searching and filtering technologies have emerged as an important enabler to assist users in navigating large amounts of content to produce and manage subsets of available content. These technologies are desired to allow users to filter content according to the user's desires or interests.\nCurrent filtering technologies can be used to generate content results according to user-defined characteristics or criteria. However, these filtering technologies may be inefficient in sufficiently allowing a user to focus the content results to the user's particular interests or goals. For example, a user may select certain genre tags to generate song results meeting the desired genre. If the results are not sufficiently narrowed to the user's desires or interests, the user may have to change or replace the genre tags to produce new results. The user may have to guess which tag changes or replacements should be made to achieve focused results with no guarantee that more focused results will be produced. The user may have to provide several iterations of different filtering settings or characteristics to finally reach a result set that is sufficiently narrowed to the user's interests. As one can imagine, this can be time consuming and possibly frustrating to a user, especially if the iterations do not provide sufficiently focused results.\nFurther, the user may not be aware of all available filter settings that can be used to filter content to the desired results. Unknown filter settings may be the ones that would successfully produce the desired focused results. For example, in the iTunes® music classification system, there are over 230 different genre and sub-genre music types. Many users are only familiar with the most popular genre types. However, the most popular types often produce the greatest number of search results and thus provide the least focus. Less commonly known genres may be available to produce the desired focused results.\nThus, a need exists for providing a user the ability to more effectively focus content results to the user's particular desires or interests. The user may be unfamiliar with all available filtering characteristics associated with the content at issue. Thus, the user may not be able to efficiently focus content results without additional suggestions or assistance."}
{"text": "Magnetoresistive random-access memory (“MRAM”) is a non-volatile memory technology that stores data through magnetic storage elements. These elements are two ferromagnetic plates or electrodes that can hold a magnetic field and are separated by a non-magnetic material, such as a non-magnetic metal or insulator. This structure is known as a magnetic tunnel junction (MTJ).\nMRAM devices can store information by changing the orientation of the magnetization of the free layer of the MTJ. In particular, based on whether the free layer is in a parallel or anti-parallel alignment relative to the reference layer, either a one or a zero can be stored in each MRAM cell. Due to the spin-polarized electron tunneling effect, the electrical resistance of the cell change due to the orientation of the magnetic fields of the two layers. The electrical resistance is typically referred to as tunnel magnetoresistance (TMR) which is a magnetoresistive effect that occurs in a MTJ. The cell's resistance will be different for the parallel and anti-parallel states and thus the cell's resistance can be used to distinguish between a one and a zero. One important feature of MRAM devices is that they are non-volatile memory devices, since they maintain the information even when the power is off.\nMRAM devices are considered as the next generation structures for a wide range of memory applications. MRAM products based on spin torque transfer switching are already making its way into large data storage devices. Spin transfer torque magnetic random access memory (STT-MRAM), or spin transfer switching, uses spin-aligned (polarized) electrons to change the magnetization orientation of the free layer in the magnetic tunnel junction. In general, electrons possess a spin, a quantized number of angular momentum intrinsic to the electron. An electrical current is generally unpolarized, e.g., it consists of 50% spin up and 50% spin down electrons. Passing a current though a magnetic layer polarizes electrons with the spin orientation corresponding to the magnetization direction of the magnetic layer (e.g., polarizer), thus produces a spin-polarized current. If a spin-polarized current is passed to the magnetic region of a free layer in the MTJ device, the electrons will transfer a portion of their spin-angular momentum to the magnetization layer to produce a torque on the magnetization of the free layer. Thus, this spin transfer torque can switch the magnetization of the free layer, which, in effect, writes either a one or a zero based on whether the free layer is in the parallel or anti-parallel states relative to the reference layer.\nFIG. 1 shows a conventional MRAM memory array architecture. Two adjacent memory cells 101 and 110 are shown. As shown in the memory cell 100, the MRAM cell has a bit line 102 and a source line 103 to write zeros and ones to an MTJ 104. As shown in cell 100, when the bit line 102 is high (e.g., Vdd) and a source line 103 is low (e.g., Vss) and a word line 106 is high, activating a gating transistor 105, current flows from the bit line 102 through the MTJ 104 to the source line 103, writing a zero in the MTJ 104. This is illustrated as the current 107. As shown in the memory cell 101, when the bit line 110 is low and the source line 111 is high and the word line 112 is high to activate the gating transistor 113, current flows from the source line 111 through the MTJ 114 (e.g., in the opposite direction) to the bit line 110, writing a one, as shown by the current 115.\nFIG. 2 shows a conventional MRAM array 200. The array 200 shows columns of cells arranged between respective source lines 240-243 and bit lines 230-233. As fabrication process sizes get smaller and smaller, more and more cells are able to be fabricated within a given die area, effectively increasing the density of a memory array. Increasing density has the benefit of more memory per unit area and less power consumption. As the cell sizes get smaller, an overriding limitation becomes the pitch width of the parallel traces of the source lines 240-243 and the bit lines 230-233. As the array becomes more and more dense with increasingly smaller fabrication processes, a limitation emerges regarding the pitch width 250-252 (e.g., the amount of distance between parallel traces) of the array. This pitch width can approach a minimum. Below the minimum jeopardizes the proper functioning of the array. This minimum pitch width can effectively halt the increasing density of memory arrays even while using increasingly smaller fabrication processes.\nThus what is needed is a way to increase densities of an MRAM array without reducing pitch width below minimums. What is needed is a way to take advantage of advancing semiconductor fabrication techniques without impinging upon the minimum pitch width limits. What is needed is a way to increase MRAM array density and thereby increase performance and reduce costs while maintaining MRAM array reliability."}
{"text": "The present invention relates to electrical connectors and, more particularly, to protective hoods for electrical connectors.\nElectrical connectors for mass connection of signals conventionally include a metal connector housing containing male or female pins through which connections are made to mating elements of a plug or jack. An insulating hood is conventionally connected to the metal housing both to provide a gripping surface for inserting and removing the connector and for strain relief for the wires exiting the connector. Strain relief is conventionally provided by either molding the mating parts of the connector to tightly fit a cable passage within the connector hood to a particular cable size or by providing a family of slip-in adapters from which a particular one can be selected to adjust the passage size to the cable size.\nOne such strain relief device is disclosed in U.S. Pat. No. 3,569,914 wherein a molded rigid planar member is held in gripping contact with the cable when the hood is assembled. The disclosed device is specific to the particular cable diameter used and a single hood is not adaptable to cable diameters which may range from those containing only two or three fine-gauge wires to those containing as many as fifty or more. Thus, to use the disclosed device on a range of cable sizes, a family of devices must be manufactured and stocked. This, of course, increases the design, manufacture and overhead cost of using such a hood.\nOne apparatus for adapting a single hood to a range of cable sizes is disclosed in U.S. Pat. No. 3,794,960 wherein a screw-driven clamp bar is urged against one side of the cable sheath to thereby capture the cable between the clamp bar and an opposed surface of a passage within the hood. This device not only requires a threaded hole in the hood, but also requires the additional parts of a clamp bar and screw. Furthermore, since the clamp bar is forcibly held in contact against a limited surface area of one side of the cable, this contact must be positioned in an area where the wires are covered by a cable sheath in order to avoid damaging the wires themselves.\nThe last-mentioned patent also illustrates a desirable feature of connector hoods; that is, a feature which permits exiting the cable from the hood either axially or laterally in order to adapt the cable routing to the user's needs. Lateral exit is enabled by providing for the optional installation of an external cable clamp for binding the cable to the hood and to thereby constrain the exit in one lateral direction.\nA further strain relief is disclosed in U.S. Pat. No. 3,966,293 wherein the wires from the connector pins are at least partly wrapped about a strain relief member which contains comb-like depressions therein helping to hold the wires. This patent also discloses exiting the wires either axially or in either or both of two lateral directions."}
{"text": "Impulse fluid or ink jets are designed and driven so as to eject a droplet of fluid such as ink from the chamber through an orifice of the ink jet device. In many applications, it is not necessary to operate the ink jet device at high performance levels, i.e., at high velocities and long throw distances. However, many applications including industrial applications require high performance ink jet devices.\nFor example, in various industrial ink jet applications, it is very important to eject droplets at high velocities with long throw distances so as to reach targets some distance from the ink jet orifice while maintaining a relatively small droplet size to create a high resolution dot on the target itself. In order to achieve this result, it is important that the head of the droplet as well as its tail remain attached to each other and travel at the same relatively high velocity.\nIn the prior art, it has been difficult to achieve high velocity and long throw distances. For example, with expanding piezoelectric transducers in ink jet print heads of the type disclosed in U.S. Pat. No.4,646,106, high performance is achieved in terms of frequency response using fluidic Helmholtz frequencies of from 25 to 50 kHz and comparable piezoelectric length mode resonant frequencies. However, the droplets formed have long tails which tend to lower the droplet velocity and the throw distance thus precluding optimum performance.\nU.S. Pat. Nos. 4,523,201 and 4,523,200 disclose similar print heads driven by voltage waveforms having a first pulse of longer duration and a second pulse of shorter duration designed to achieve early break off of the droplet tail. However, the devices disclosed therein are designed to operate at Helmholtz frequencies of less than 50 kHz and there is no disclosure of the effect of exciting higher harmonic frequencies to achieve break off of the tail for producing higher velocity droplets with improved throw distance. Rather the second pulse merely improves aiming.\nReference is now made to FIGS. 1A through 1F which schematically depict the drive waveform in FIG. 1A and the ink jet device itself at various points in time in FIGS. 1B through 1F. Referring to FIG. 1A, the drive waveforms depicted with voltage on the ordinate and time on the abscissa. At time A, the ink jet device as depicted in FIG. 1B is maintained in the quiescent state with the transducer 10 unenergized and a predetermined volume of ink 12 contained within the chamber 14 behind an orifice 16. At time B as shown in FIG. 1A, the transducer 10 is driven by the voltage pulse 17 as shown so as to contract the length of the transducer 10 thereby increasing the volume of ink 12 within the chamber 14 and pulling back the meniscus 18 in the orifice 16 to the position shown.\nAs shown in FIG. 1C, which corresponds to time C of FIG. 1A, the transducer 10 begins to expand as the voltage is reduced as applied to the transducer 10. As a result, the volume of ink 12 within the chamber 14 begins to contract while advancing the meniscus 18 through the orifice 16 as shown in FIG. 1D. At a slightly later time than C but before D as shown in FIG. 1E, the transducer 10 has nearly returned to the quiescent state as shown in FIG. 1E and a droplet 20 with a ligament 21 has begun to form at the orifice 16. In FIG. 1F which corresponds to time D in FIG. 1A, the droplet 20 has traveled some distance from the orifice 16 with a slowly moving tail 22 attached. As depicted in FIG. 1F, the tail 22 has just broken off from the meniscus 18 at the orifice 16 before the volume of ink within the chamber 14 returns to the condition shown in FIG. 1B. As can be readily seen from FIG. 1F, the tail 22 is elongated in a manner so as to create a \"lay over\" condition on a target assuming the tail 22 and the head 20 remain attached throughout their flight to the target. Tail 22 which is relatively slow moving as compared to the head velocity which makes the tail grow in length and break up thereby decreasing the overall throw distance to the target.\nIn the devices of the prior art of the type disclosed in U.S. Pat. Nos. 4,459,601, 4,509,059, 4,646,106 and 4,697,193, ink jet devices have been characterized by Helmholtz resonant frequencies of approximately 40 kHz and piezo length mode resonant frequencies of 45 kHz. The tail which was formed at the meniscus saw pressure disturbances of approximately 45 kHz. As a consequence, the tail would be broken off as shown in FIG. 1F in response to this disturbance during the negative velocity part of the cycle so as to provide a very low acceleration component thereby producing a drop with a high head velocity and low tail velocity and a tail that grows longer over long print gaps producing poor print quality."}
{"text": "This invention relates to an apparatus for assembling a door for a vehicle, and is more particularly directed to an assembling apparatus having a set jig for positioning and setting thereon a door provided to the side of an assembling station. The set jig is arranged to be movable to advance toward and retreat from a door opening portion of a vehicle body set on the assembling station so that the door may be assembled with the door opening portion by an advance movement of the set jig.\nUsually with this kind of apparatus, as disclosed, for instance, in Japanese Unexamined Patent Application Publication No. Sho 54-5276, a set jig is supported in a floating condition on a frame base which is movable to advance toward and retreat from a vehicle body. A locating member which is to be brought into engagement with a reference opening made in the vehicle body is fixedly provided on the set jig so that before the advance movement of the frame base, the locating member is brought into engagement with the reference opening. While so engaged, the set jig may be moved with the movement of the locating member so as to be set in position at its position corresponding to the door opening portion.\nSuch a construction, however, is inconvenient in that, because the locating member is forcibly inserted into the reference opening and causes the set jig to move therewith, the vehicle body side is liable to be injured. It often happens that an excessively large force is applied to the locating member. In addition, if a balance mechanism for lightening the movement of the set jig is provided, the apparatus becomes complicated in construction and high in price."}
{"text": "Technical Field\nThe present invention relates to a windshield glass peripheral structure.\nRelated Art\nIn vehicles such as automobiles, a windshield glass and a cowl louver positioned at a vehicle from side of the windshield glass are coupled together with a molding interposed therebetween. Japanese Patent Application Laid-Open (JP-A) No. 2015-51653 describes a coupling structure between a windshield glass and a cowl louver that is provided with a molding that includes a fixed portion fixed to an inner face of a lower edge portion of the windshield glass, and a fitting portion that receives an anchor rib that protrudes downward from a location near to an upper edge portion of the cowl louver.\nOrdinarily, a molding provided between a windshield glass and a cowl louver is joined to an outer face of the windshield glass. In contrast thereto, in the invention described in JP-A No. 2015-51653, an interposed portion of the molding is interposed between the outer face of the windshield glass and an outer face of the cowl louver; and the outer face of the windshield glass and the outer face of the cowl louver, and a leading end of the interposed portion of the molding, are disposed so as to be flush with each other. In a case of an ordinary molding, when water from a vehicle width direction end portion of the windshield glass enters at the inside of the windshield glass, water droplets run along an inner face of the windshield glass and flow down toward the vehicle up-down direction lower side, falling as they are toward the cowl.\nHowever, in case of the invention described in JP-A No. 2015-51653, the fixed portion of the molding is joined to the inner face of the windshield glass and the molding is curved protruding toward the vehicle front and vehicle bottom. Thus, some of the water droplets that have run along the inner face oldie windshield glass and flowed down toward the vehicle up-down direction, lower side at the vehicle width direction end portion of the windshield glass then run along an upper end portion of the fixed portion of the molding and head toward the vehicle width direction inside. It is therefore possible that the water droplets drip onto part of a wiper unit, such as a wiper link, provided further to the vehicle inside than the cowl louver. There is therefore room for improvement to suppress early deterioration of the wiper unit due to water exposure of the wiper link.\nIn order to address the problem above, conceivably, for example, a water stopping member such as a dam rubber might be provided to the upper end portion of the molding, to dam water droplets from running along the upper end portion of the molding and beading toward the vehicle width direction inside. However, in such a configuration, providing the water stopping member to the molding as a separate body increases the number of components, thereby increasing the assembly steps and component costs.\nIn consideration of the above, an object of the present invention is to obtain a windshield glass peripheral structure capable of suppressing water exposure of a wiper unit due to water droplets falling from a molding without adding any new components."}
{"text": "While global efforts for environmental protection of the Earth progress, strengthening of regulations regarding flue gas or fuel efficiency in internal combustion engines, such as engines of automobiles is continuing. Since turbochargers can reduce engine emissions compared to natural aspiration by sending compressed air into an engine, turbochargers are devices that are very effective in fuel efficiency improvement and CO2 reduction.\nIn turbochargers, a turbine is rotationally driven by engine flue gas, thereby rotating a delivery fan of a coaxial centrifugal compressor. The air compressed by the rotation of the delivery fan is raised in pressure by being reduced in speed by a diffuser, and is supplied to the engine through a scroll flow path. In addition, regarding methods for driving turbochargers, not only driving using flue gas, but also various forms, such as driving using an electric motor and driving a prime mover via a transmission, are known.\nIf flow rate decreases in the centrifugal compressor used for such a turbocharger, surging that is the pulsation of the whole system occurs. In order to expand the operating range of a centrifugal compressor, it is necessary to reduce a critical flow rate at which surging occurs.\nIn contrast, PTL 1 discloses a technique of forming a recirculation flow path that allows a portion of the air sucked into a delivery fan to circulate therethrough, thereby expanding the operating range. This recirculation flow path is formed in a compressor cover that covers the delivery fan. This recirculation flow path is formed when an inner tube member having fins disposed at intervals in a peripheral direction is externally fitted to a compressor cover body in an axial direction. In this case, when outer peripheral ends of the fins abut against the compressor cover from the inside, the recirculation flow path is partitioned into a plurality of sections in the peripheral direction. A portion of the air sucked into the delivery fan is introduced into the recirculation flow path as a swirling flow, and is straightened by the above fins."}
{"text": "1. Field of the Invention\nThe present invention relates to signal processing, and, in particular, to computer-implemented processes and apparatuses for encoding and decoding image signals for progressive transmission and display.\n2. Description of the Related Art\nStill images and video images typically require large numbers of bits to represent digitally, even using sophisticated compression techniques. The time required to transmit such images for display at a remote destination, and therefore the time delay between display of successive images, may prove disturbing to the remote viewer. Using a conventional progressive transmission technique, such as those based on the wavelet transform, may prove computationally intensive and therefore time consuming.\nIt is desirable to provide encoding systems for generating, encoding, and transmitting image signals and decoding system for receiving, decoding, and displaying image signals that reduce the delay to the remote viewer. In particular, it is desirable to provide personal computer (PC) based conferencing systems that provide the capabilities for efficient transmission of images from one conference participant to a remote conference participant over relatively low bandwidth media, such as a PSTN telephone line.\nIt is, therefore, an object of the present invention to provide computer-implemented processes and apparatuses for efficiently generating, encoding, and transmitting image signals and methods, apparatuses, and systems for efficiently receiving, decoding, and displaying image signals.\nIt is a particular object that the present invention be applicable to PC-based conferencing systems.\nFurther objects and advantages of this invention will become apparent from the detailed description of a preferred embodiment which follows.\nThe present invention is a computer-implemented process and apparatus for encoding image signals. According to a preferred embodiment, a location within an image is selected and signals corresponding to the image are encoded following a spatial decomposition pattern, wherein the spatial decomposition pattern is based on the selected location.\nThe present invention is also a computer-implemented process and apparatus for displaying encoded image signals corresponding to an original image. According to a preferred embodiment, an initial set of encoded image signals is provided and a display image is displayed in accordance with the initial set of encoded image signals. A sequence of subsequent sets of encoded image signals is provided and the display image is progressively updated in accordance with the sequence, wherein the sequence corresponds to a spatial decomposition pattern based on a selected location within the original image.\nThe present invention is also a computer-implemented process and apparatus for transmitting image signals corresponding to an original image at a local node for display at a remote node. According to a preferred embodiment, a location within the original image is selected at the local node. The image is divided into a plurality of blocks at the local node and a transform is applied to each of the blocks to generate a plurality of transformed blocks at the local node, wherein each of the transformed blocks comprises a DC transformed signal and a plurality of AC transformed signals. The DC transformed signals for all of the transformed blocks are encoded to generate encoded DC transformed signals at the local node and the encoded DC transformed signals are transmitted from the local node to the remote node. The encoded DC transformed signals are decoded at the remote node to generate decoded DC transformed signals and a display image is displayed in accordance with the decoded DC transformed signals on a monitor at the remote node. The AC transformed signals for all of the transformed blocks are encoded following a spatial decomposition pattern at the local node to generate a sequence of encoded AC transformed signals, wherein the spatial decomposition pattern is based on the selected location, and the sequence of encoded AC transformed signals is transmitted from the local node to the remote node. The sequence of encoded AC transformed signals is decoded at the remote node to generate a sequence of decoded AC transformed signals and the display image is progressively updated in accordance with the sequence of decoded AC transformed signals."}
{"text": "Water-dispersible pressure sensitive adhesive compositions have many uses. For example, such adhesives are useful with garments such as gowns, sheets, drapes, and other mammalian body coverings, etc. which are required during surgery and other aseptic and non-aseptic procedures. Many of these procedures are carried out in the normal operation of hospitals, medical and dental clinics and the like. Additionally, such adhesives are useful with sterilization indicator tapes.\nIn such applications, it is essential that the adhesive be resistant to body fluids normally encountered during use, yet be dispersible under alkaline conditions such as are commonly encountered during laundering.\nA number of water dispersible adhesive systems are known. For example, U.S. Pat. No. 5,125,995 (D'Haese et al), WO 93/06184 (D'Haese), U.S. Pat. No. 4,413,080 (Blake), U.S. Pat. No. 4,569,960 (Blake), and copending application U.S. Ser. No. 07/889,647 filed May 28, 1992 (Patnode et al).\nExamples of other water-dispersible pressure sensitive adhesives include U.S. Pat. No. 3,865,770 (Blake), U.S. Pat. No. 3,441,430 (Peterson) and U.S. Pat. No. 2,838,421 (Sohl).\nA variety of other pressure-sensitive adhesives are said to be either water-soluble or water-dispersible. See, for example U.S. Pat. No. 4,413,082 (Gleichenhagen et al), U.S. Pat. No. 4,341,680 (Hauber et al), European Patent Publication 0 352 442, U.S. Pat. No. 3,152,940 (Abel et al), U.S. Pat. No. 4,338,432 (Eskay), U.S. Pat. No. 3,096,202 (deGroot von Arx), European Patent Publication 0 297 451 (Knutson et al), U.S. Pat. No. 3,763,117 (McKenna et al), and U.S. Pat. No. 3,890,292 (Bohme et al).\nThe present invention provides a pressure sensitive adhesive composition and articles that employ the adhesive which can be used in a variety of applications. The adhesive is especially well suited for preparing pressure sensitive adhesives tapes which can be used with products such as gowns, sheets, drapes, etc. used in hospitals and medical and dental clinics. The adhesives can also be used in a variety of other applications such as sterilization indicator tapes and labels."}
{"text": "The present invention relates to equipment for the testing of printed circuit boards (PCB's) for electronic circuits.\nMore specifically, the invention pertains to printed circuit test equipment of the type having spring loaded contact fingers distributed on a contact-bearing plate and having a contact selecting plate which, according to the type of circuit to be tested, lets predetermined ones of the contact fingers protrude, whereby the contact fingers can be set to the particular printed circuit board.\nVarious types of equipment for testing printed circuit boards by establishing an electric contact between certain contact lands and contact fingers which can be connected with test equipment are known from DE-AS 26 57 910, DE-AS 22 43 457, DE-OS 33 11 977, CH-PS 633 403. In these types of equipment, contact-bearing plates are used which, due to their permanent arrangement of the contact fingers, cannot be used except for a single type of printed circuit board, or as in DE-AS 22 43 457 require expensive adaptations for other types of printed circuit boards. The expense for the production of test-specific substitution parts for the adaptation of the equipment is relatively high. This is a prticular disadvantage if a small number of printed circuit boards are to be produced and tested.\nA further disadvantage of the above-mentioned conventional equipment is that separate wiring is necessary for each of the contact fingers, which are arranged on the contact-bearing plate according to the requirements of the type of PCB to be tested, with the test equipment or an adapter which connects the test equipment with the PCB. Even utilizing the so-called wire-wrap-technique, such wiring is expensive. Furthermore, large numbers of substitution sets are required for a certain type of printed circuit boards and such sets may in some instances have to be stored for decades for various applications.\nRecently equipment of the above-mentioned type has become known, in which the contact fingers are arranged on the contact bearing plate in a pre-determined distribution, i.e. in a regular matrix arrangement with a standardized distance (module). The selection of the contact fingers necessary for the testing of a specific type of printed circuit board is effected via a contact selecting plate which is equipped with fixed selecting fingers. These selecting fingers are also arranged according to the basic matrix, but only at the locations of those contact fingers which are necessary for the respective test. In the assembled state of contact bearing plate and contact selecting plate, the contact fingers necessary for the testing are released from a rest position to contact with a printed circuit board to which the arrangement is set by producing underpressure. The contacting of the contact fingers used for the testing of the printed circuit board with the test equipment or the interface respectively, etc., in these conventional arrangements is effected by the wiring of the respective selecting fingers, which again results in the disadvantages already mentioned.\nThus, the number of substitution parts of the equipment necessary for each kind of printed circuit boards to be tested is reduced, and moreover, the substitution parts are less complicated as well. In addition to the above-mentioned disadvantage of wiring necessity the problems concerning storage of the substitution parts are also similar to those of other equipment already mentioned, because the contact-selecting plate with the selecting fingers--in order to prevent damage--either has to be packed accordingly or stored at a proper distance to the other substitution parts.\nA further disadvantage is that the protruding contact fingers at the contact selecting plate make it much more difficult to adapt the equipment automatically for the testing of another type of printed circuit boards or in practice even make it impossible to adapt the equipment accordingly, thus requiring expensive manual adaptation."}
{"text": "This relates to the protection of integrated circuits from electrostatic discharge (ESD). More particularly, it relates to the protection of differential circuits from ESD.\nESD protection has been a main concern in the reliability of integrated circuit products in various sub-micron technologies. ESD is the transient discharge of static charge that can arise from activities such as human handling, machine contact or field-induced charging of a packaged IC. Specific models have been developed to represent these discharges such as the Human Body Model (HBM), the Machine Model (MM), and the Charged Device Model (CDM), respectively. See, for example, A. Amerasekera and C. Duvvury, ESD in Silicon Integrated Circuits, pp. 17-40 (2d Ed., Wiley, 2002), which is incorporated herein by reference.\nFIG. 1 is a schematic diagram of an illustrative differential circuit 100 with conventional ESD protection circuitry. Differential circuit 100 comprises first and second transistors 110, 130, a clamp transistor 150 and a diode 170. Illustratively, first and second transistors 110, 130 are a low voltage differential signaling (LVDS) output pair. Each transistor 110, 130, 150 is a MOS transistor with a source and drain formed in a body of the transistor and an insulated gate over the body in the region between the source and drain. In the schematic diagram of FIG. 1, the bodies of transistors 110, 130, and 150 are identified as elements 112, 132, 152; the sources are identified as elements 114, 134, 154; the drains are identified as elements 116, 136, 156; and the gates are identified as elements 118, 138, 158, respectively. Sources 114 and 134 are connected together at a source node 190. Resistors 120, 140, 160 are schematic representations of the circuitry between gates 118, 138, 158, respectively, and a common node 180; and resistor 182 is a schematic representation of the circuitry between source node 190 and common node 180. As is known in the art, the actual circuitry represented by these resistors may be considerably more complicated than a simple resistance. Input terminals 122, 142 are connected to gates 118, 138, respectively; and output terminals 124, 144 are connected to drains 116, 136, respectively. Diode 170 may be implemented as a dedicated diode or as the body diode of a MOSFET clamp transistor similar to transistor 150 or as both devices connected in parallel.\nAs is known in the art, the differential circuit typically comprises several other circuit elements not shown in FIG. 1. For example, transistor 110 is typically driven by circuitry connected to input terminal 122. Additionally, other circuits are connected to output terminal 124 to pull this node up when transistor 110 is in the off state.\nTypically, the transistors of differential circuit 100 are NMOS transistors with a P-type body and N-type source and drain regions. As a result, since the P-type body and the N-type source region of each transistor form a first P-N junction and the P-type body and the N-type drain region form a second P-N junction, a parasitic lateral bipolar transistor is present in each transistor. In the event of a positive voltage ESD event on the output terminal 124, circuit 100 is intended to operate so that the second P-N junction of clamp transistor 150 is driven into breakdown and avalanche and the parasitic transistor is triggered into conduction to discharge the ESD pulse.\nHowever, during the ESD event, the body voltage of transistors 110 and 130 can easily float above the source voltage, also making possible bipolar triggering of transistors 110 and 130. For example, as shown in the voltage vs. time plot of FIG. 2A, in the case of a positive ESD event on output terminal 124, the voltages on nodes 180 and 190 will both rise until the voltage on node 180 reaches the threshold voltage of transistor 130 at time t1. Transistor 130 then begins to pull down the voltage at node 190 while the voltage at node 180 is basically pinned at one Vbe above ground by diode 170. While the voltage on node 190 keeps decreasing, the avalanche current in transistor 110 and the voltage at output terminal 124 keep increasing. Eventually, destructive bipolar triggering will occur in transistor 110 when the body-source junction becomes fully forward biased at time t2 leading to a rapid drop in the output voltage. To prevent this, the clamp transistor 150 must trigger before it happens; but it is difficult to assure consistent, timely triggering without significant additional circuitry."}
{"text": "This invention relates to an apparatus for the manufacture of a flexible heating appliance, comprising at least one heating wire constituting an electrical heating circuit being placed between textile materials which are combinable with each other thermally, directly or with the aid of an additional adhesive layer, and combining the said textile materials by the application of heat and pressure.\nFlexible heating appliances of the aforesaid kind have been known, for instance, from German Pat. No. 1,757,215 which provides for inserting heating conductors into woven-in channels in a textile fabric .In part, the heating conductors are already woven-in or are received between weldable, channel-forming layers, while the welding process itself is carried out by special high-frequency welding tools devised for this type of manufacturing method.\nA simple method of fixation comprises the use of coated fabrics, webs or fleeces which are combinable in flat shape with each other, in the manner used in making outer wear or over-clothes. It is furthermore known to provide grooves or similar recesses in an ironing plate or ironing board, into which grooves the heating conductor being placed on the underneath layer, is pressed during the ironing treatment, so that there is no combination such as by welding in the zone of the groove or the like recess, and the heating conductor remains displaceable in the channel system being formed. Within the program of a further development of flexible heating bodies which can be used for heating seats, and in particular automobile seats, it was necessary to achieve a short heating-up period in order to provide a satisfactory heating of the seat surface, if possible within a very short time. A prerequisite for attaining this goal is an inherent softness of the flexible heating body which latter must lie directly under the cover material of the seat, while the heating circuit design must not become contrasted on the surface of the cover material.\nFor this purpose, thin heating conductors insulated with coatings having the properties of perfluoropolyethylenes such as Teflon, and whose diameter is smaller than 2 mm, have been sewn on to a lower layer of material, and have then been additionally affixed in position by covering layers of material being joined to the layer underneath by ironing. The layout in an ironing plate mentioned hereinbefore can be carried out much more rapidly, but it is necessary to this end that the nails which determine the layout of the heat carrier and its bends are depressable in grooves to be preferably flush with them during the ironing treatment.\nBoth of these known methods are expensive, as a lack of fit between the ironing plate and the nail template or stencil may cause serious damage to the template."}
{"text": "The present invention relates to a heat engine construction. More particularly, the present invention relates to an engine construction including a heat exchanger which is movable between two extreme positions within an enclosed space, and which produces a variation in pressure which may be transmitted to a working piston.\nVarious types of engines are known which are employed to operate as heat engines, with such engines being described, for example, in U.S. Pat. No. 3,956,894 to Tibbs; U.S. Pat. No. 3,878,680 to Dauvergne; U.S. Pat. No. 4,077,221 to Maeda; and U.S. Pat. No. 4,270,351 to Kuhns.\nBy the present invention, there is provided an internal thermal exchanger engine which, in one embodiment, includes an enclosed shell or chamber having hot and cold end portions and having contained therein a heat exchanger which occupies approximately one-half the total volume of the chamber, the heat exchanger being open ended and containing heat conductive materials such as fine copper wire strands which function in heat exchange relationship with a gas contained in the chamber. In an alternative embodiment of the invention, the enclosed chamber is in the form of a semi-cylinder, with the heat exchanger of a wedge-shape construction and being freely rotatable between a hot side of the semi-cylinder and a cold side across the arc of the semi-cylinder.\nIn the present heat engine, the pressure differential is achieved in a most efficient way due to the action of the internal heat exchanger which removes energy from the gas as the exchanger moves from cold to hot end locations, storing such energy within the system in the exchanger material, creating low pressure in the chamber and then releasing energy to the gas to create high pressure in the chamber as the exchanger moves from the hot to the cold end locations within the chamber.\nIt is an object of the present invention to provide a thermal engine which is capable of highly efficient operation.\nIt is a further object of the invention to utilize to the fullest extent the energy supplied during operation of a heat engine.\nIt is another object of the invention to provide a thermal engine which operates within an enclosed space to produce a pressure differential which may be transmitted to a location outside the enclosed space.\nIt is another object of the invention to re-cycle heat energy within the heat engine system by the use of an internal heat exchanger.\nIt is another object of the invention to provide a system which removes energy from a gas contained within an enclosed chamber and which transfers this energy to a working piston."}
{"text": "The present invention relates to an apparatus for protectively mounting an object. Such protection is afforded when the object rests on a shock-mounted support that is isolated from the ground by means of resilient shock absorbing elements. Shock mountings can be designed to absorb or to attenuate, to an acceptable value, redundant forces and dampen vibrations that would otherwise be transmitted into the object and possibly damage it. These redundant forces could arise from shock or impact due to severe or careless handling, by vibration such as could be encountered when the object is being carried on a moving conveyance, or by combinations of shock and vibration. The design of shock absorbers is accomplished by applying well known equations of vibration dynamics to the parameters of the particular system of components and elements involved. In the present disclosure, the term shock absorbers will be understood to include resilient members capable of absorbing vibrations as well as shock.\nA known technique for protecting machinery, delicate instruments, and the like is to place them onto or in a shock absorbing support apparatus or container that has been designed for the application. Although the particular object is protected after it is securely installed into its protective means, the very procedure of installing the object can itself cause damage, particularly if the object is heavy and/or delicate. Manual, or external mechanical lifting devices such as fork lifts, cranes, ramps, pry bars and the like are usually employed, and the possibility of imposing a force or forces severe enough to damage the still unprotected object is apparent. Similarly, during the procedure of removing the object from its shock absorbing apparatus it is again vulnerable to damage, perhaps more so since loading equipment available at the location where the object was prepared for shipping often times is not available at the location where the object is to be installed."}
{"text": "1. Field\nApparatuses and methods consistent with exemplary embodiments relate to a user terminal device and a displaying method thereof, and more particularly, to a user terminal a user terminal device including a display panel in a square shape and a displaying method thereof.\n2. Description of the Related Art\nWith the development of electronic technology, various types of user terminal devices have been developed and distributed. Recently, a size of user terminal devices has been minimized, but the functions of the user terminal devices have become more diverse, and, therefore, the demands for the user terminal devices have been increasing.\nA user terminal device may provide various content, such as multimedia contents and application screens, upon a request of a user. A user may select a desired function using a button or a touch screen provided on a user terminal device. A user terminal device may selectively execute a program according to a user interaction between a user and the device, and display the result thereof.\nAs more diverse functions are provided by a user terminal device, there is a need for a method of displaying content or a user interface method. That is, as the method of displaying content has changed and the type and function of the content have been increasing, the related art interaction methods such as selecting a button or touching a screen may be insufficient to perform various functions of a user terminal device.\nThus, there is a need of user interaction technology which enables a user to use a user terminal device more conveniently."}
{"text": "1. Field of the Invention\nThe present invention relates to a device for locking the jaws of an adjustable wrench around a pipe to minimize the possibility that the wrench will slip or fall loose from the pipe. More particularly, the invention relates to a detachable jaw-locking device for a pipe wrench having an extendable locking member which extends across the jaw opening of the wrench to enclose the open end of the wrench and lock the jaws around the pipe.\n2. The Prior Art\nThe use of pipe wrenches at substantial heights by workmen on a construction site is a common occurrence. The fact that such wrenches have an open side and rely on applied pressure from the user to maintain a grip on the pipe, creates serious safety hazzards to both workmen and other persons below the work site. Often, the wrenches are carelessly released or may simply be left engaged around a pipe or other object. In either case, the wrenches can easily slip, or otherwise become dislodged and fall. This risk is further exaggerated by the substantial weight of the pipe wrench as compared with other wrenches or tools of common use.\nU.S. Pat. Nos. 504,232 and 2,517,041 disclose wrenches having spring biased jaws for better gripping, but neither was intended to solve the additional problems described hereabove. Even though spring biased jaws help lessen the risk of dislodging a pipe wrench engaged around a pipe, the jaw opening of the engaged wrench provides a potential escape route for any such wrench when released, jarred or otherwise disturbed.\nFrom the foregoing, it will be appreciated that what is needed in the art is a device which is capable of enclosing the open end of an adjustable pipe wrench and firmly locking its jaws around a pipe, regardless of the size of the jaw spread. Such an invention is described herebelow."}
{"text": "The present invention provides a simplified and reliable construction for a high-pressure rotating water jet nozzle which is particularly well suited to industrial uses where the operating parameters can be in the range of 1,000 to 40,000 psi, rotating speeds of 1000 rpm or more and flow rates of 2 to 50 gpm. Under such use the size, construction, cost, durability and ease of maintenance for such devices present many problems. Combined length and diameter of such devices may not exceed a few inches. The more extreme operating parameters and great reduction in size compound the problems. Pressure, temperature and wear factors affect durability and ease of maintenance and attendant cost, inconvenience and safety in use of such devices. Use of small metal parts and poor quality of materials in such devices may result in their deterioration or breakage and related malfunctioning and jamming of small spray discharge orifices or the like. The present invention addresses these issues by providing a simplified construction with a greatly reduced number of parts and a design in which net operating forces on nozzle components are minimized."}
{"text": "A traditional skateboard fundamentally consists of a board, normally made of wood to which four wheels in pairs are fixed, which serves for practicing the sport called skateboarding. The wheels are grouped in two units of two wheels, the axis of each one being held with a flexible rod slightly inclined in relation to the board which allows it to carry out turns by inclining the board from one side to another and the smooth turning movement of the axes which it entails.\nThe U.S. Pat. No. 6,056,302 from Marc Smith describes an example of a wheel unit for a traditional skateboard. As can be observed in the figures, the unit is formed by a base fixed to the board and this unit to the axis of the wheels by means of a rod which passes through bushings. This configuration allows for the occurrence of a slight inclination of the wheels in relation to the board when the user is supported on one or the other side of said board which allows smooth turns to be carried out in the direction of displacement.\nRecently, a new type of skateboard has emerged called the “surfskate”. The difference between a traditional skateboard and a surfskate is fundamentally the type of front wheel unit used. The front wheel unit of a surfskate is configured to allow much greater maneuverability than a traditional skateboard and it is also equipped with a centering spring which exerts an action on the wheels intended to return them to the neutral position thereof. A surfskate has very different behavior to that of a conventional skateboard, allowing the user to propel themselves by simply carrying out oscillatory hip movements similar to those carried out when surfing.\nThe U.S. Pat. No. 6,793,224 from Carver Skateboard shows an example of a wheel unit for a surfskate. This unit comprises a base couplable to the lower part of a board and an arm coupled to the base which rotates in relation to the same around a first axis. A rod which has a pair of wheels mounted at opposing ends is fixed to said arm and can be rotated in relation to the same around a second axis. A compression spring connected between the base and the arm limits the rotational movement of the arm and returns it towards a central position aligned to the direction of the moving surfskate.\nAlthough the use of this type of wheel unit significantly increases the maneuverability of a surfskate with respect to that of a traditional skateboard, the rotation thereof continues being somewhat insufficient. In addition, it must be separated from the board by means of separators enough in order to avoid the wheel impacting the board. Ultimately, the construction of the unit is complicated and this makes the maintenance and repair operations difficult to carry out."}
{"text": "1. Field of the Invention\nThe present invention relates to safety devices, specifically to a helmet having a safety buckle system.\n2. Description of the Related Art\nA helmet is a form of protective gear worn on the head to protect it from injuries. All helmets attempt to protect the user's head by absorbing mechanical energy and protecting against penetration. Their structure and protective capacity are altered in high-energy impacts. Beside their energy-absorption capability, their volume and weight are also important issues, since higher volume and weight increase the injury risk for the user's head and neck. Anatomical helmets adapted to the inner head structure were invented by neurosurgeons at the end of the 20th century.\nHelmets used for different purposes have different designs. For example, a bicycle helmet must protect against blunt impact forces from the wearer's head striking the road. A helmet designed for rock climbing must protect against heavy impact, and against objects such as small rocks and climbing equipment falling from above. Practical concerns also dictate helmet design: a bicycling helmet should be aerodynamic in shape and well ventilated, while a rock climbing helmet must be lightweight and small so that it does not interfere with climbing.\nIn addition, a football helmet is a protective device used primarily in American football and Canadian football. It consists of a hard plastic top with thick padding on the inside, a face mask made of one or more rubber coated metal bars, and a chinstrap. Some players add polycarbonate visors to their helmets, which are used to protect eyes from glare and impacts. Helmets are a requirement at all levels of organized football, except for non-tackle variations such as flag football. Although they are protective, players can and do still suffer head injuries such as a concussion. Each position has a different type of face mask to balance protection and visibility.\nSome improvements have been made in the field. Examples of references related to the present invention are described below in their own words, and the supporting teachings of each reference are incorporated by reference herein:\nU.S. Pat. No. 6,003,156, issued to Anderson, discloses a Safety Helmet Lockout System comprised of a safety helmet with a uniquely keyed system contained within its cranial webbing and chinstrap that has the ability to sense when said helmet is securely strapped upon the wearers head. The chinstrap contains a locking buckle which releases a uniquely keyed enabling key and/or a signal upon sensing that said safety helmet is securely strapped upon the wearers head. Several means of sensing are possible, but the preferred embodiment employs physio sensors located within the cranial webbing and chinstrap of said safety helmet. The locking buckle which is located on the chinstrap, releases said enabling key or signal after receiving signals from all the physio sensors, and sensing tension in the chinstrap. The released enabling key, or emitted signal, can then be used to access vehicles, sporting goods, industrial equipment, tools, passage locks, and the like. The access locks can range from simple ignition locks and passage locks, to retracting pin or keyed brake systems, keyed clutch systems, and specialty lockout systems for certain applications (such as, but not limited to, skiis, snowboards, and snowmobiles).\nU.S. Pat. No. 5,946,735, issued to Bayes, discloses a quick-release football helmet chin strap that incorporates the use of a hook and loop fastener to tighten the strap and secure the helmet to the player's head. Intended for use on newly produced helmets and replacing the snap-type fasteners on existing helmets, the use of the hook and loop fasteners eliminates the burden associated with using the snap fasteners during the frequent between-play attachment and detachment.\nU.S. Pat. No. 4,279,037, issued to Morgan, discloses an adjustable suspension for protective headgear wherein a constantly uniform tension on all securing members of the suspension can be simultaneously increased or decreased uniformly through a one-step tension control device without removing the headgear from the wearer's head. A chin-strapless embodiment which anchors against vertical displacement on the forehead and skull base. Detachable chin strap with concealed dual locking fastener. A headgear wherein the protective elements are supported by an adjustable, detachable suspension.\nU.S. Patent Application Publication No.: 2010/0319701, by Connell, discloses a respirator assembly has a respirator mask and a separate bracing device. The bracing device has a pressure element mountable over a frontpiece of the respirator mask, which bears against that frontpiece to hold the respirator mask against the face of a user. The bracing device also has two connection portions, each having one end attached to the pressure element and the other end being for attaching the bracing device to a helmet worn by the user. The connection portions are extendable. Thus, shortening or distending the connection portions reduces the distance between the helmet and the respirator mask and thereby causes the pressure element to bear against the frontpiece of the respirator mask, compressing the user's head between the helmet and the mask and thereby holding the mask against the user's face.\nThe inventions heretofore known suffer from a number of disadvantages which include being unsafe, being unable to securely couple a helmet to a wearer, being difficult to use, being uncomfortable, being unable to securely couple a chin strap to a helmet; being easily removed; being easily knocked off; coming unsnapped on impact; coming uncoupled when pressure is applied on the strap in a forward direction; being difficult to detach; and the like.\nWhat is needed is a helmet that solves one or more of the problems described herein and/or one or more problems that may come to the attention of one skilled in the art upon becoming familiar with this specification."}
{"text": "A single conventional order for goods or services may comprise the purchase of a variety of different kinds of goods or services. Several different suppliers may be required to fulfil a particular order. Accordingly; several different carriers may be required to transport the goods from supplier through to the final delivery address or addresses. The final delivery addresses might be the same as or different than an installation address. In some situations, no delivery or installation is required."}
{"text": "The present invention is generally related to medical devices, kits, and methods. More specifically, the present invention provides a system for crossing stenosis, partial occlusions, or total occlusions in a patient's body.\nCardiovascular disease frequently arises from the accumulation of atheromatous material on the inner walls of vascular lumens, particularly arterial lumens of the coronary and other vasculature, resulting in a condition known as atherosclerosis. Atheromatous and other vascular deposits restrict blood flow and can cause ischemia which, in acute cases, can result in myocardial infarction or a heart attack. Atheromatous deposits can have widely varying properties, with some deposits being relatively soft and others being fibrous and/or calcified. In the latter case, the deposits are frequently referred to as plaque. Atherosclerosis occurs naturally as a result of aging, but may also be aggravated by factors such as diet, hypertension, heredity, vascular injury, and the like.\nAtherosclerosis can be treated in a variety of ways, including drugs, bypass surgery, and a variety of catheter-based approaches which rely on intravascular widening or removal of the atheromatous or other material occluding the blood vessel. Particular catheter-based interventions include angioplasty, atherectomy, laser ablation, stenting, and the like. For the most part, the catheters used for these interventions must be introduced over a guidewire, and the guidewire must be placed across the lesion prior to catheter placement. Initial guidewire placement, however, can be difficult or impossible in tortuous regions of the vasculature. Moreover, it can be equally difficult if the lesion is total or near total, i.e. the lesion occludes the blood vessel lumen to such an extent that the guidewire cannot be advanced across.\nTo overcome this difficulty, forward-cutting atherectomy catheters have been proposed. Such catheters usually can have a forwardly disposed blade (U.S. Pat. No. 4,926,858) or rotating burr (U.S. Pat. No. 4,445,509). While effective in some cases, these catheter systems, even with a separate guidewire, have great difficulty in traversing through the small and tortuous body lumens of the patients and reaching the target site.\nFor these reasons, it is desired to provide devices, kits, and methods which can access small, tortuous regions of the vasculature and which can remove atheromatous, thrombotic, and other occluding materials from within blood vessels. In particular, it is desired to provide atherectomy systems which can pass through partial occlusions, total occlusions, stenosis, and be able to macerate blood clots or thrombotic material. It is further desirable that the atherectomy system have the ability to infuse and aspirate fluids before, during, or after crossing the lesion. At least some of these objectives will be met by the devices and methods of the present invention described hereinafter and in the claims."}
{"text": "Traditional RF tuners may be used to receive terrestrial or cable broadcast signals in the range of about 50 megahertz (MHz) to 860 MHz. Such tuners may employ single conversion intermediate frequency (IF) technology. A single conversion tuner may mix a received RF signal with a local oscillator (LO) signal in a single mixing stage to produce an IF signal. A common IF frequency may be about 36 MHz, for example.\nThe tuner may include a number of tracking filters used to protect a desired channel from interfering signals. Certain interfering signals, including an image channel and intermodulation products, may be generated in a subsequent stage without the tracking filters. The image channel may lie above the desired channel at twice the IF output frequency, for a tuner utilizing high side mixing. The image channel may be particularly problematic, because it may lie directly on a desired channel after an IF conversion. A tracking filter may be centered on the desired channel to suppress one or more undesired channels prior to conversion.\nSome tuner front-ends may include RF circuitry for a number of sub-bands, for coverage of a full broadcast spectrum. An RF circuit associated with each sub-band may include a tracking filter of substantially identical design to tracking filters associated with the other sub-bands.\nSome tuners may include a first tracking filter comprising a single inductor-capacitor resonant network located prior to an automatic gain control (AGC) stage. A center frequency associated with the first tracking filter may track with a center frequency of the desired channel. The first tracking filter may select the desired channel from the full received spectrum and may provide a first attenuation to the undesired channels, including the image channel. This may provide protection from intermodulation that might otherwise be generated in a subsequent stage.\nA second tracking filter may comprise two resonant networks arranged as a double-tuned, loosely-coupled structure. The second tracking filter may be located prior to a conversion stage. A center frequency associated with the second tracking filter may track the desired channel frequency. The second tracking filter may provide a higher resonant quality factor (Q) attenuation to the undesired channels, including the image channel.\nIn a typical implementation, the tracking filters may be replicated for each segmented band. The tracking filters may be formed from numerous discrete components, including varactor diodes and air coils. An inductance value associated with the air coils may vary from a few hundred nanohenries (nHs) in the low band to a few nHs in the high band. The air coil inductors may be characterised by a high Q factor, typically in excess of 50, and may require manual adjustment. A tracking filter thus composed may not be suitable for semiconductor integration."}
{"text": "The present invention relates to residential heating and/or cooling systems and other indoor comfort systems, and is more particularly concerned with battery-powered thermostats of the type that derive the power for monitoring conditions within the comfort space and for controlling the signaling to the furnace or other comfort system from a battery i.e, dry cells or alkaline cells. The invention is more specifically directed to DC thermostats with latching relays that are pulsed to change their state from OFF to ON or from ON to OFF, and remain latched into that state until pulsed into the other state. Heat may be provided from a gas, oil, or electric furnace or heat pump, and cooling may be provided from a compression/condensation/expansion/evaporation cycle air conditioner, an absorption type air conditioner, a ground-water heat exchange cooing system, or other available chilling apparatus. As used here, the term “cooling” includes both sensible cooling (reducing the temperature of the comfort air) and latent cooling (removing humidity). These comfort air apparatus may have additional functions for better control of the environmental air in the comfort space, such as multiple fan speeds, high and low heating, and high and low compressor speeds.\nWall thermostats are typically installed on an interior wall of a dwelling, business space or other residential or commercial space to control the operation of a furnace, air conditioner, heat pump, or other environmental control equipment. The thermostat continuously monitors the temperature of the room or other interior comfort space or zone, and is connected by a run of thermostat wires to the associated environmental control equipment to signal a call for heating, a call for cooling, or otherwise to keep an interior comfort space parameter, such as temperature, within some range (e.g., 68° F., ±1.0° F.). Other controls may be available, sensitive to other parameters, such as humidity or particulate level.\nIn the examples discussed here, the thermostats are powered by small dry-cell batteries, e.g., alkaline power cells. However, the re-pulsing action can be used favorably with any other thermostat type as well, including those powered by 24 volt AC thermostat power or any other energy source.\nIn many applications, i.e., in many permanent homes, and in mobile, recreational, or marine dwellings and spaces, battery-powered thermostats may be the preferred thermostat. These typically have a battery power source, e.g., a pair of AA alkaline power cells, installed within the thermostat housing to power the electronics. These supply DC energy to the internal electronics within the thermostat and power latching relays that connect thermostat power (such as 24 volts AC) to the particular thermostat wires that control heating, fan, air conditioning, and so forth. The thermostat is constantly monitoring the temperature in the comfort space and is also constantly monitoring the voltage level available on the DC battery power source. One or more controls, i.e., push buttons, rotary knob(s) etc., allow the occupant to set and adjust temperature setpoints for heating and/or for air conditioning. An LCD display screen on the housing of the thermostat allows the occupant to see the room temperature and also to see other functions, such as temperature setpoints during a temperature setting sequence. A low-battery message can also be displayed on this screen if the battery voltage drops below a level that indicates the battery is approaching the end of its useful life and should be replaced.\nThe latch mechanism in most latching relays relies on the magnetic force of the contact to keep the contact latched ON or latched OFF. If the thermostat is bumped, the mechanical shock can sometimes move the contact from the desired state. In that case, the microprocessor in the thermostat continues as if the relay were still in its desired state, but the furnace or air conditioning equipment does not. For example if a low temperature induces the thermostat to issue a call for heat, the microprocessor pulses the latching relay associated with heat to its ON condition, and the relay contact applies AC thermostat power to the W or heat thermostat wire, causing the furnace to commence a heat cycle. Then, if a bump or other shock knocks the contact back to the OFF condition, the furnace will shut down as if the thermostat had been satisfied. Wrong position of the relay could also occur due to a sudden or momentary power supply drop, component degradation, or other reasons that may cause the relay to fail to move to the intended position. Meanwhile, the thermostat heat set point has not been reached, and the microprocessor does not act on the heating relay. With the furnace shut down, the air in the comfort zone continues to cool down, but the thermostat does not respond because the thermostat's call for heat has not been satisfied.\nOn the other hand, thermal runaway can result if the heat relay contact is bumped into its ON condition when there is no call for heat, as the thermostat microprocessor acts as if the heat relay were open of OFF.\nIt may be the case that a the above-described erroneous condition will occur when the occupant is away for an extended period, and that the comfort space will continue to cool down until the occupant returns and can correct the thermostat. Failure of the furnace or other HVAC apparatus to cycle properly during winter conditions can result in frozen water pipes or other damage, including temperature stress to indoor plants and pets\nBecause of the foregoing problems, it would be desirable to provide a feature or features for battery powered thermostats that will ensure the relays are properly latched into the proper conditions, and that the heating and/or cooling conditions in the comfort space not be permitted to vary significantly during a heating or cooling operation."}
{"text": "1. Field of the Invention\nThe present invention relates to the transfer molding of semiconductor devices. More specifically, the present invention relates to a method of using a dam in transfer molding encapsulation of a semiconductor device and a heat sink, and in the resulting semiconductor device assembly.\n2. State of the Art\nA semiconductor integrated circuit (IC) device (referred to as a die or chip) includes bond pads on the active surface thereof for interfacing the integrated circuits of the semiconductor device with other circuits outside the die located on differing substrates. Since the semiconductor devices are relatively small and the attendant bond pads on the active surface thereof are, in comparison, considerably smaller, lead frames having a plurality of leads thereon connected to the bond pads of a semiconductor device are used to connect the semiconductor device with other circuits on differing substrates.\nIn a conventional lead frame design for use with an integrated circuit semiconductor device, the lead frame includes a plurality of leads having their ends terminating adjacent a side or edge of the integrated circuit semiconductor device with the device being supported by the die paddle portion of the lead frame. Electrical connections are made by means of wire bonds extending between the leads of the lead frame and the bond pads located on the active surface of the integrated circuit semiconductor device.\nSubsequent to the wire bonding operation, portions of the leads of the lead frame and the integrated circuit semiconductor device may be encapsulated in suitable plastic material to form a packaged semiconductor device assembly. The leads and lead frame are then trimmed and formed to the desired configuration after the packaging of the semiconductor device in the encapsulant material.\nIn a Leads-Over-Chip (LOC) type lead frame configuration for an integrated circuit semiconductor (IC) device assembly, the leads of the lead frame extend over the active surface of the semiconductor device being insulated therefrom by tape which is adhesively bonded to the semiconductor device and the leads of the lead frame. Electrical connections are made between the leads of the lead frame and bond pads on the active surface of the semiconductor device by way of wire bonds extending therebetween. After wire bonding, the leads of the LOC lead frame and the semiconductor device are encapsulated in suitable plastic to encapsulate the semiconductor device and portions of the leads. Subsequently, the leads are trimmed and formed to the desired configuration to complete the packaged semiconductor device.\nBy far the most common manner of forming a plastic package about a semiconductor device assembly is molding and, more specifically, transfer molding. In this process, with specific reference to a LOC type semiconductor die assembly, a semiconductor die is suspended by its active surface from the underside of inner lead extensions of a lead frame (typically Cu or Alloy 42) by a tape, screen print or spin-on dielectric adhesive layer. The bond pads of the die and the inner lead ends of the frame are then electrically connected by wire bonds (typically Au, although Al and other metal alloy wires have also been employed) by means known in the art. The resulting LOC die assembly, which may comprise the framework of a dual-in-line package (DIP), zig-zag in-line package (ZIP), small outline j-lead package (SOJ), quad flat pack (QFP), plastic leaded chip carrier (PLCC), surface mount device (SMD) or other plastic package configuration known in the art, is placed in a mold cavity and encapsulated in a thermosetting polymer which, when heated, reacts irreversibly to form a highly cross-linked matrix no longer capable of being re-melted.\nThe thermosetting polymer generally is comprised of three major components: an epoxy resin, a hardener (including accelerators), and a filler material. Other additives such as flame retardants, mold release agents and colorants are also employed in relatively small amounts. While many variations of the three major components are known in the art, the focus of the present invention resides in the filler materials employed and their effects on the active die surface.\nFiller materials are usually a form of fused silica, although other materials such as calcium carbonates, calcium silicates, talc, mica and clays have been employed for less rigorous applications. Powdered, fused quartz is currently the primary filler used in encapsulants. Fillers provide a number of advantages in comparison to unfilled encapsulants. For example, fillers reinforce the polymer and thus provide additional package strength, enhance thermal conductivity of the package, provide enhanced resistance to thermal shock, and greatly reduce the cost of the polymer in comparison to its unfilled state. Fillers also beneficially reduce the coefficient of thermal expansion (CTE) of the composite material by about fifty percent in comparison to the unfilled polymer, resulting in a CTE much closer to that of the silicon or gallium arsenide die. Filler materials, however, also present some recognized disadvantages, including increasing the stiffness of the plastic package, as well as the moisture permeability of the package.\nWhen a heat sink is used on a semiconductor device assembly package, encapsulation of the semiconductor device becomes more difficult during the transfer molding process. In the first instance, the inclusion of the heat sink along with the semiconductor device attached to the lead frame makes the transfer molding of the assembly more difficult as more components must be placed and aligned within the mold cavity. Misalignment of the semiconductor device and the heat sink within the mold cavity may cause bleeding and flashing of the resin molding compound over the heat sink. Furthermore, when the heat sink, which is usually copper or an alloy thereof, rests against the mold surface during the transfer molding process, damage to the mold surface can occur by the mold surface being scratched and/or worn from contact therewith by the heat sink. The resulting worn mold surfaces cause the resin molding compound to bleed and flash over the outside of the heat sink during the transfer molding process. This affects the ability of the heat sink to transfer heat to the surrounding environment during the operation of the semiconductor device as well as presenting a poor appearance of the molded semiconductor device assembly.\nAccordingly, there is a need for an improved transfer molding process for packaging semiconductor devices having heat sinks associated therewith to help prevent or reduce the bleeding or flashing of the molding compound over portions of the heat sink during the transfer molding process of the semiconductor device assembly.\nThe present invention is directed to a semiconductor device assembly that includes a heat sink adjacent to a die. A dam positioned about the peripheral edges of the heat sink during the transfer molding process serves to help prevent damage to the mold and help prevent encapsulant packaging compound material from flowing onto the heat sink. The dam may be a resilient, non-metallic material. The dam may also be a protrusion created by, for example, stamping the heat sink from a sheet of material or stamping the bottom of the heat sink to form the dam thereon. The dam may include a suitable metal material, such as copper, copper alloys, etc., and a suitable non-metallic material, such as polyamides, and tape. The dam may be permanently connected to the heat sink or removed following packaging. The dam may be removed with heat or during an electrolytic deflash cycle or, if desired, mechanically. The invention may be employed in connection with various types of lead frame configurations or, when a lead frame is not used with a bare semiconductor device of the semiconductor device assembly."}
{"text": "1. Field of Invention\nThe present invention relates to a disposable beverage cup, and more particularly to a disposable lid for the disposable beverage cup, which provides a spill proof ability to prevent any beverage leak especially when the disposable beverage cup is tipped.\n2. Description of Related Arts\nA conventional disposable lid of a disposable beverage cup on the current market, which is usually made of plastic, typically has a drinking hole arranged for a user to drink the beverage in the disposable beverage cup, and a vent hole that allows air venting into the disposable beverage cup therethrough for maintaining the balance of the pressure between the inside and outside of the disposable beverage cup. However, the conventional lid has several drawbacks.\nA major drawback of the conventional lid is that the conventional lid provides a relatively weak sealing and leak proof ability. In other words, the disposable beverage cup with the conventional lid is inconvenient for carrying. Most of lids in the current market have exposed drinking holes and vent holes. An improved lid may further comprise a valve only provided at the drinking hole. However, the structure of the valve is complicated which will highly increase the manufacturing cost of the lid. In fact, none of the existing lid provides any seal to seal the vent hole of the lid.\nSince the size of the drinking hole is larger than the size of vent hole, the beverage in the disposable beverage cup will be spilled out through the drinking hole when the disposable beverage cup is tipped. Accordingly, the beverage in the disposable beverage cup may be spilled onto a table, a floor, clothes, file documents, or inside a car, resulting in unnecessary trouble and loss. Worse yet, if the beverage is hot, consumer may get seriously burnt.\nFurthermore, since the consumer must pay more attention on carefully holding the disposable beverage cup, he or she may loss focus on other things, which may cause potential hazard. For example, when the consumer carefully holds the disposable beverage cup during driving, he or she may loss focus on driving which may cause accident. Since the beverage is easy to leak out, some people may even pour out the beverage onto a road, or discard the disposable beverage cup onto a green area aside a road or a lawn of a front yard of a house and polluting our environment.\nAnother solution found in some fast-food restaurants or beverage shops is to provide a beverage tray which is made of cardboard for retaining the disposable beverage cup at the upright position. A beverage tray can typically hold four to six disposable beverage cups. It will be a waste when the customer only buys one or two disposable beverage cups but the seller still has to supply an oversized tray. It will increase the cost and produce additional trashes and possibly environment pollution."}
{"text": "The present invention relates to a filter medium for filtering impurities out of fluids and a method of manufacturing the filter medium.\nA filter medium for filtering impurities having relatively large diameters is generally formed by a kind of sieve. However, in a fluid there are contained not only solid grains having large diameters, but also particles having very small diameters and fluids of different kinds. For instance, in an engine oil for use in motor cars there are included various kinds of particles having different sizes such as iron powders, carbon particles, pigments and liquids such as water. Particularly, fine particles and liquids could not be easily removed by the conventional filter medium. Therefore, absorbing agents such as active carbon and activated white clay(kaoline) are provided in the filter medium.\nHeretofore, there has been proposed a charcoal filter in which absorbing agents are mixed in a filtering member or inserted between sheet-like filtering members such as filtering paper and nonwoven fabric. In these known filter medium, it is quite difficult to perform the desired filtering action due to the fact that the absorbing agents might move in the filtering members or flow out of the filtering members.\nIt has been further known that the hot active carbon particles are spread on an nonwoven fabric made of polypropylene resin and are fused onto the fabric. However, in such a filter medium, the absorbing agents, i.e. carbon particles are not firmly secured to the filtering member i.e. polypropylene unwover fabric and further the filterability of the fabric is largely lost at positions where the carbon particles are fused, so that the filtering efficiency is decreased. It has been also known to secure the absorbing agents to the filtering member by means of a bonding agent. In such a known filter medium, the filtering efficiency is also decreased. Further, the bonding agent might affect the absorbing function of the absorbing agents.\nIn an oil cleaner, fine magnetic particals such as iron particles produced by the friction of machine parts should be removed. Such magnetic particles may be absorbed by means of a filter medium including permanent magnet particles. However, it is almost impossible to spread uniformly the magnet particles or powder on the filtering member, because the magnet particles are attracted to each other.\nFurther, in case of filtering impurities out of an oil, it is necessary to remove water contained in the oil. However, in the known filter medium since the filtering member is clogged due to the water absorption by means of OH-radical, the filterability is decreased to a large extent."}
{"text": "Merchants are always looking for new channels through which to market their products to potential customers. Marketing channels such as newspapers; magazines; television; radio; posters; signs on bus-stop benches; signs on taxi cabs, buses, and other forms of public transportation; and so-called “junk mail” have been in wide use for many years. One marketing channel that has come along more recently is the portion of the Internet known as the World Wide Web (the “Web”).\nThe Internet also supports a wide variety of electronic-commerce systems, including secure wireless-payment systems that permit a user to pay for goods and services using a mobile device such as a cellular telephone. Though such mobile electronic-commerce systems are increasing in popularity, such systems have significant untapped potential as channels through which marketing messages can be transmitted to mobile-device users."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor element typified by a field effect transistor to be formed over a substrate having a low strain point and to a method for manufacturing the same, and relates to a semiconductor integrated circuit including the semiconductor element and to a method for manufacturing the same. Specifically, the present invention relates to a thin film transistor in which a gate insulating film is heat-treated at a temperature beyond a strain point of a substrate such as glass and to a method for manufacturing the same.\n2. Description of Related Art\nIn recent years, development of a system-on-panel incorporating a logic circuit such as a CPU or a memory as well as a pixel or a driver circuit over a light transmitting insulating substrate such as glass or quartz has been attracting attention. High-speed operation is required for a driver circuit and a logic circuit, and manufacturing a thin film transistor (hereinafter, also referred to as a TFT) having high switching speed is required in order to realize it. It is effective for realizing a TFT having higher switching speed to use a semiconductor film with fewer crystal defects as an active layer; to make a gate insulating film thinner, and to miniaturize a transistor size typified by miniaturization of a gate length.\nCharacteristics required for a gate insulating film can be given as follows: few defects in a thin film; without a fixed charge; a low interface level with a semiconductor film; low leakage current; and the like. However, a gate leakage current tends to increase with a decrease in a film thickness of a gate insulating film. In addition, such a fine gate insulating film that can lower a gate leakage current is required in order to make the gate insulating film thinner. A field effect semiconductor device that can be driven with a low voltage and responds well to a high drive frequency can be obtained by making the gate insulating film thinner.    Patent Document 1: Japanese Patent Laid-Open No. H6-188421"}
{"text": "Handheld showerheads typically have showerhead and handle portions. The showerhead portion includes a showerhead face with nozzles and openings for delivering water to a user from the handheld showerhead. The handle portion provides a structure for a user to hold when using the handheld showerhead.\nHandheld showerheads may include more than one mode of operation. Multiple modes of operation provide a user with flexibility to select a desired spray pattern, or pause water flow from the handheld showerhead. Some possible spray patterns for a handheld showerhead with multiple modes of operation may include standard water streams, converging water streams, pulsating water streams, and mist sprays. For a handheld showerhead with multiple modes of operation, a circular ring is formed to rotate around the showerhead face. A user rotates the circular ring around the showerhead face until the desired mode of operation is selected."}
{"text": "Mercury can be present in trace amounts in all types of hydrocarbon streams such as crude oils. The amount can range from below the analytical detection limit to several thousand ppbw (parts per billion by weight) depending on the source.\nMethods have been disclosed to remove mercury from liquid hydrocarbon feed, specifically volatile mercury. It has been reported that mercury in crude is primarily in the form of volatile species, e.g., 90% Hg0 and only 10% DMHg. See Wilhelm et al. Energy & Fuels 2006, 20, 180-186 (See Table 5 on page 184) http://pubs.acs.org/doi/pdf/10.1021/ef0501391. In contrast, Tao et al. J. Anal. At. Spectrom., 1998, 13, 1085-1093 http://pubs.rsc.org/en/content/articlepdf/1998/ja/a803369b show in Table 8 that dialkylmercury species out-proportion elemental Hg in condensates and natural gas liquids.\nU.S. Pat. No. 4,915,818 discloses a method of removing mercury from liquid hydrocarbons (natural gas condensate) by contact with a dilute aqueous solution of alkali metal sulfide salt, for reaction of the sulfur component with the mercury, where the mercury sulfur compounds precipitate and settle for subsequent recovery as a solid waste. U.S. Pat. No. 6,268,543 discloses a method for removing elemental mercury with a sulfur compound where mercury is removed as a solid. U.S. Pat. No. 6,350,372 discloses the removal of mercury from a hydrocarbon feed upon contact with an oil soluble or oil miscible sulfur compound U.S. Pat. No. 4,474,896 discloses using polysulfide based absorbents to remove elemental mercury (Hg0) from gaseous and liquid hydrocarbon streams.\nThere are also a number of commercially available processes and products for the removal of (volatile) elemental mercury Hg0 from hydrocarbon streams including but not limited to ICI Synetix' Merespec™ fixed bed absorbents, UOP's HgSIV™ regenerative mercury removal adsorbents, and Johnson Matthey's Puraspec™ and Puracare™ granulated absorbents for the removal of mercury from naphtha and/or gaseous hydrocarbon streams. Adsorption technology does not work well for crude oils and condensates with low levels of mercury, and with primarily non-volatile mercury.\nProduction of oil and gas is usually accompanied by the production of water. This produced water in some cases is reinjected into the subsurface for disposal, to maintain pressure in the reservoir, or to achieve other beneficial effects. The produced water may consist of formation water (water present naturally in the reservoir), or water previously injected into the formation. As exploited reservoirs mature, the quantity of water produced increases. Produced water is the largest single fluid stream in exploration and production operations.\nThere is a need for methods for the removal of non-volatile mercury from liquid hydrocarbon streams, and particularly methods wherein produced water can be used/recycled."}
{"text": "1. Field of the Invention\nThe present invention relates to computer security and electronic mail. More specifically, the present invention relates to a method and an apparatus for facilitating transmission of an encrypted electronic mail message to anonymous recipients without divulging the identities of the anonymous recipients.\n2. Related Art\nThe advent of computer networks has led to an explosion in the development of applications that facilitate rapid dissemination of information. In particular, electronic mail (email) is becoming the predominant method for communicating textual and other non-voice information. Using email, it is just as easy to send a message to a recipient on another continent as it is to send a message to a recipient within the same building. Furthermore, an email message typically takes only minutes to arrive, instead of the days it takes for conventional mail to snake its way along roads and through airports.\nOne problem with email is that it is hard to ensure that sensitive information sent through email is kept confidential. This is because an email message can potentially traverse many different computer networks and many different computer systems before it arrives at its ultimate destination. An adversary can potentially intercept an email message at any of these intermediate points along the way.\nOne way to remedy this problem is to “encrypt” sensitive data using an encryption key so that only someone who possesses a corresponding decryption key can decrypt the message. (Note that for commonly used symmetric encryption mechanisms the encryption key and the decryption key are the same key.) A person sending sensitive data through email can encrypt the sensitive data using the encryption key before it is sent through email. At the other end, the recipient of the email can use the corresponding decryption key to decrypt the sensitive information.\nEncryption works well for a message sent to a single recipient. However, encryption becomes more complicated for a message sent to multiple recipients. This is because encryption keys must be managed between the sender and the multiple recipients.\nConventional mail protocols, such as the Pretty Good Privacy (PGP) protocol, send mail to multiple recipients by encrypting a message with a session key (that is randomly selected for the message) to form an encrypted message. The session key is then encrypted with the public key of each of the recipients to form a set of encrypted keys. This set of encrypted keys is sent with the encrypted message to all of the recipients. Each recipient uses one of its private keys to decrypt an encrypted session key and then uses the session key to decrypt the encrypted message.\nNote that key identifiers for the public keys that were used to encrypt the encrypted session keys are sent along with the encrypted session keys, so that each recipient can determine whether or not the recipient possesses a corresponding private key that can decrypt the encrypted session key. These identifiers are typically generated by computing as a hash of the public key.\nUnfortunately, the key identifiers can also identify a recipient of an email message to other recipients of the email message. This complicates the process of sending an encrypted email message to anonymous recipients, because the recipients of the email message can determine the identities of the anonymous recipients by examining the key identifiers for the anonymous recipients.\nWhat is needed is a method and an apparatus for facilitating transmission of encrypted email to anonymous recipients without divulging the identities of the anonymous recipients."}
{"text": "Zeolites are crystalline or quasi-crystalline aluminosilicates constructed of repeating SiO4 and AlO4 tetrahedral units. These units are linked together to form frameworks having regular intra-crystalline cavities and channels of molecular dimensions. Numerous types of synthetic zeolites have been synthesized and each has a unique framework based on the specific arrangement of its tetrahedral units. By convention, each framework type is assigned a unique three-letter code (e.g., “CHA”) by the International Zeolite Association (IZA).\nSynthetic CHA zeolites are produced using a structure directing agent (SDA), also referred to as a “template” or “templating agent”. SDAs are typically complex organic molecules which guide or direct the molecular shape and pattern of the zeolite's framework. Generally, the SDA serves to position hydrated silica and alumina and/or as a mold around which the zeolite crystals form. After the crystals are formed, the SDA is removed from the interior structure of the crystals, leaving a molecularly porous aluminosilicate cage.\nZeolites have numerous industrial applications including internal combustion engines, gas turbines, coal-fired power plants, and the like. In one example, nitrogen oxides (NOx) in the exhaust gas may be controlled through a so-called selective catalytic reduction (SCR) process whereby NOx compounds in the exhaust gas are contacted with a reducing agent in the presence of a zeolite catalyst.\nZSM-5 and Beta zeolites have been studied as SCR catalysts due to their relatively wide temperature activity window. However, the relatively large pore structures of these zeolites have a number of drawbacks. First, they are susceptible to high temperature hydrothermal degradation resulting in a loss of activity. Also, large and medium pore sizes tend to adsorb hydrocarbons which are oxidized as the temperature of the catalyst increases, thus generating a significant exotherm which can thermally damage the catalyst. This problem is particularly acute in lean-burn systems, such as vehicular diesel engines, where significant quantities of hydrocarbon can be adsorbed during cold-start. Coking by hydrocarbons presents another significant drawback of these relatively large and medium pore molecular sieve catalysts. In contrast, small pore molecular sieve materials, such as those having a CHA framework type code (as defined by the International Zeolite Association), offer an improvement in that fewer hydrocarbons are able to permeate into the framework.\nTo promote the catalytic reaction, transition metals may be included in the zeolite material, either as a substituted framework metal (commonly referred to as “metal-substituted zeolite”) or as a post-synthesis ion exchanged or impregnated metal (commonly referred to as “metal-exchanged zeolite”). As used herein, the term “post-synthesis” means subsequent to zeolite crystallization. The typical process for incorporating a transition metal into a zeolite is by cationic exchange or impregnation of metals or precursors after the molecular sieve is formed. However, these exchange and impregnation processes for incorporating metals frequently lead to poor uniformity of metal distribution and the smaller pores of CHA type molecular sieve materials exacerbate that problem.\nSo-called “one pot” synthesis procedures in which a transition metal compound is present during synthesis of the molecular sieve framework vis-à-vis present as a post-synthesis ion exchanged metal. However, known one-pot processes lack sufficient control over metal loading, yield framework structures with inadequate silicon to aluminum ratios (SAR), and/or necessarily include alkali metal in the synthesis mixture which can poison acid sites and have a detrimental effect on the hydrothermal stability. Moreover, reported one-pot synthesis procedures for forming metal-containing molecular sieves have been observed to yield significant amounts, in some instances as much as 20%, of amorphous phase, copper oxides and other impurities, which negatively impact the stability and activity of the catalyst."}
{"text": "The present invention relates generally to centrifuge systems and, more particularly, to an advanced design self-compensating centrifuge arm in a centrifuge system for testing precision instrumentation devices."}
{"text": "One example of conventional decorative sheets is such that a protective layer formed of an ionizing radiation-cured resin, such as an electron beam-cured resin, is provided. In this type of decorative sheets, a pattern, such as a woodgrain pattern, is provided on the surface of a substrate sheet, such as paper, an olefin resin, or a vinyl chloride resin, and a protective layer formed of an electron beam-cured resin or the like is provided on the surface of the pattern. In general, such decorative sheets are applied, for example, as interior materials for fittings or buildings, to particle boards, plywoods, plastic members, or other substrates for decorative plates, and then used as decorative plates.\nIn conventional decorative sheets provided with a protective layer formed of an ionizing radiation-cured resin, the ionizing radiation-cured resin has a regulated crosslink molecular weight. In general, molecular weight between crosslinks or average crosslink molecular weight is used as an index of the crosslink molecular weight. For conventional electron beam-cured resins, for example, the average molecular weight between crosslinks is set to 150 to 1,000, more preferably 200 to 1,000, still more preferably 250 to 800. In this case, the average molecular weight between crosslinks is determined by the following equation:Average molecular weight between crosslinks=molecular weight of whole resin (m)/number of crosslink points \nIn this equation, the molecular weight of the whole resin is Σ (number of moles of each component incorporated x molecular weight of each component), and the number of crosslink points is Σ[2(number of functional groups in each component −1) x number of moles of each component].\nWhen the average molecular weight between crosslinks exceeds 1,000, the whole resin is excessively soft. In this case, high hardness characteristic of electron beam-cured resins is lost, posing a problem of scratch resistance of the protective layer in its surface. Even though the average molecular weight between crosslinks is 250 to 800, which has been regarded as a preferred molecular weight range for the conventional protective layer, for example, a protective layer having an average molecular weight between crosslinks of about 800 close to the upper limit value poses a problem of resistance to staining. In this case, in a test on resistance to staining, a contaminant is likely to be left on the surface of the protective layer. On the other hand, when the average molecular weight between crosslinks is brought to not more than 200, the resistance of the protective layer to staining is significantly improved. The reason for this is that the hardness of the protective layer in its surface has been increased by virtue of an increase in proportion of the functional group to the molecular weight. This enables the amount of the contaminant left on the surface to be reduced. When the average molecular weight between crosslinks is not more than 200, the electron beam-curable resin coheres at the time of curing, unfavorably leading to a failure of the print layer of a woodgrain pattern or the like underlying the protective layer. Consequently, a failure of adhesion occurs. In particular, when the ink composition constituting the print layer contains a large amount of a pigment, the separation between layer is likely to occur.\nThe present invention has been made with a view to solving the above problems of the prior art, and it is an object of the present invention to provide a decorative material which has a protective layer possessing excellent resistance to staining and does not cause a problem of a failure of adhesion or the like."}
{"text": "The technical field of this invention is vehicle power steering systems.\nBasic power steering systems for vehicles provide a fixed steering assist offset curve regardless of driving conditions. More sophisticated power steering systems, however, are available that are capable of changing power steering effort in different driving situations. For example, one such system, marketed by Delphi Automotive Systems, provides a decrease in steering assist with vehicle speed, so that greater assist will be provided at very low speed vehicle maneuvers such as parking, while significantly less assist is provided at high speed highway driving for a tighter, xe2x80x9cmanual steering feelxe2x80x9d and greater stability. In a recent version of this system, a variation responsive to estimated lateral acceleration has been added to improve steering linearity.\nAnother such system is described in U.S. Pat. No. 6,062,336, issued May 16, 2000 and entitled Adaptive Variable Effort Steering System. The power steering system described therein takes advantage of the advent of vehicle low friction road surface handling controls for vehicle brakes and/or suspension components to provide special control of power steering assist near or at the limit of vehicle handling. Such controls include anti-lock braking controls (ABS), traction controls (TCS) and vehicle stability enhancement (VSE) controls. Each of these controls responds to sensors indicating when a vehicle approaches or reaches a vehicle handling limit, where traction is decreased or lost, to override normal vehicle control while the normal handling is not optimal. The power steering system of U.S. Pat. No. 6,062,336 is responsive to signals from one or more such controls that indicate when the control has become active in modifying handling, and thus indicate the existence of a near limit or at limit vehicle handling situation. When such a handling limit signal is received in the system of the patent, the power steering system responds by decreasing steering assist, and thus increasing steering effort, to provide a more xe2x80x9cmanualxe2x80x9d steering feel as long as the handling limit situation is indicated.\nBut vehicles with an adaptive power steering system as described do not experience loss of traction only on low coefficient road surfaces. They may also lose traction if driven in a xe2x80x9cperformancexe2x80x9d mode (in aggressive or evasive maneuvers) on high coefficient surfaces. But on a high coefficient surface, high lateral forces can be generated on the tires and transmitted through the steering system to produce an inherent increase in steering effort; and an additional increase in steering effort can produce an unpleasant sensation to a driver, particularly at the initiation of a counter-steering input. Thus, the increase in steering effort beneficial in improving driver comfort near or at the handling limit on low coefficient of friction surfaces has proven to be unnecessary and undesirable in similar circumstances on high coefficient surfaces. To modify the added steering downward for the system\"\"s behavior on high coefficient surfaces can reduce its desired effect on low coefficient surfaces.\nThe system described in U.S. Pat. No. 6,062,336 is described as having embodiments responsive to variable vehicle dynamic parameters such as yaw rate error, in the case of a vehicle stability enhancement system, to provide a variable magnitude signal of approaching vehicle handling limit; but this signal does not distinguish between loss of traction on a low coefficient surface and loss of traction on a high coefficient surface.\nA variable power steering system for a motor vehicle is responsive to an activity signal from a vehicle stability enhancement (VSE) system and a lateral surface coefficient of friction estimator to modify steering assist to provide a greater steering assist when the VSE system is active on a low coefficient surface than when the VSE system is active on a high coefficient surface. Thus, the power steering assist is automatically varied, when the vehicle stability enhancement system is active, to provide the optimal handling feel on both low coefficient and high coefficient surfaces.\nPreferably the modification is phased in and out in slew limited steps that also decrease in size as lateral surface coefficient increases. In a preferred embodiment, a power steering assist correction comprises the sum of a vehicle speed offset value varying inversely with vehicle speed and a road surface adaptation value that is the greatest in magnitude of a plurality of offset values comprising a vehicle stability enhancement offset value determined as described above, a lateral acceleration offset value varying inversely with estimated lateral acceleration and offset values used during the activity of any other vehicle traction limit handling controls such as anti-lock braking (ABS) or traction control (TCS)."}
{"text": "The present invention is generally related to hypodermic syringes and more particularly to syringes that include a needle that is retractable after the intended use to substantially prevent inadvertent exposure to the needle and reuse of the syringe.\nHypodermic syringes are widely used in the medical arts for administering medicaments and for drawing body fluid samples. Generally, hypodermic syringes have a metal needle attached either fixedly or removably that has a sharpened distal point for penetrating vial stoppers or patient\"\"s skin. The hypodermic syringes and needles have been used for many years with few problems reported when the vast numbers and needles being used are considered. More recently, with the recognition of viral diseases that are transmitted by body fluids and greater sensitivity of the need to protect health care workers from inadvertent contact with previously used needles (commonly referred to as xe2x80x9csharpsxe2x80x9d) as well as the need to reduce criminal misuse of improperly disposed of needles and syringes, syringes and needles that include provisions to prevent reuse have been developed.\nProvisions intended to prevent reuse of needles and syringes include a variety of sharps collector systems that are widely used in health care facilities. Other developments include needle attachments that may be readily broken off by practitioners once the syringe has completed its intended use. A variety of shielding mechanisms has been developed; some of which are currently commercially available. While many of these developments have reduced the incidence of inadvertent exposure of healthcare workers to sharps, most of these devices can readily be overcome by an individual determined to obtain and misuse a hypodermic syringe and needle. As a result of this problem, further developments in the art of hypodermic syringes have resulted in syringes with needles that withdraw into the body of the syringe once their intended use is completed.\nU.S. Pat. No. 4,838,869 discloses a retractable hypodermic needle configured for one time use wherein the needle is spring loaded and automatically irretrievably retracted into the hypodermic syringe when the syringe plunger is fully depressed, whereby protrusions on the end of the plunger engage tabs holding the spring loaded needle to release the needle for retraction. A potential problem with the design disclosed in this patent is that many times a practitioner may draw and expel a fluid several times during preparation for administration of a medicament, with this design, the practitioner could inadvertently discharge the retraction mechanism. Further, the design would be very difficult to manufacture in large volumes.\nU.S. Pat. No. 4,900,307 discloses a hypodermic needle with an enlarged hub that provides provisions for selectively withdrawing the needle into the hub once the syringe and needle have completed their intended usage. While this disclosed design does substantially eliminate the problem of premature discharge of the retraction mechanism, the enlarged hub has a considerable xe2x80x9cdead volumexe2x80x9d that would result in a significant undeliverable retention of the medicament. Additionally, although the needle is secured in the hub after discharge, the syringe itself is still fully functional after the hub with the needle inside is removed.\nU.S. Pat. No. 4,994,034 discloses a hypodermic injection system with a retractable needle wherein the needle retracts within the interior cavity of a syringe plunger. The disclosed invention includes a cylindrical spring housing with resilient fingers which capture a coiled spring that biasly holds a needle holder against the retaining force of the resilient fingers. The plunger in this disclosure has a frangible end, which when engaging the resilient fingers under a pre-determined amount of force, dissociate which remaining inwardly-tapered shoulders spread the resilient fingers, allowing the coiled spring to eject the needle and its holder into the interior cavity of the syringe plunger. A syringe manufactured using this disclosure would be complex and difficult to assemble. It is believed that no successful commercial product has been produced using this disclosure.\nU.S. Pat. No. 5,019,044 discloses a safety hypodermic syringe with a hypodermic needle fixed connected to a holder plate and constantly supported by a spring for making axial movement. The holder plate is normally retained by a clamp at a ready position for injection. When the plunger of the syringe is pushed to the bottom of the barrel, the needle is released from the clamp and is pushed by the spring to drop and further follow a rubber plug to be squeezed into a chamber in the plunger. Again, no successful commercial product has resulted from this disclosure, which would be complex to manufacture and appears to have a considerable undeliverable dead volume.\nAnother example of a syringe with a retractable needle is disclosed in U.S. Pat. No. 5,053,010. The disclosed syringe retracts the needle into a hollow plunger additional pressure on the plunger after the contents of the syringe are expelled. The disclosed design incorporates a sliding elastomeric seal which displaces from its forward position to a retracted position, thereby allowing additional forward travel of the plunger to actuate the retraction mechanism. A problem reported with this design is that, because of the soft nature of the seal, the seal may be prematurely displaced during its use in an injection. Attempts to overcome this difficulty by increasing the stiffness of the sealing member could impair the seal integrity.\nU.S. Pat. No. 5,180,369 discloses a self-destructive syringe assembly having a needle cannula fixed to a slidable piston. The slidable piston and slidable piston flange are held within the barrel of the syringe assembly by a compressed spring, a guide tube and a shatter ring. The plunger of the syringe assembly is a hollow elongated tube with a thumb flat at one end, a sliding gasket, a plunger shatter plate and a hook rim at the other end. The patent reports that when medicament is injected, the elongated hollow plunger is further thrust into the shatter ring, the shatter ring shatters, further allowing the slidable piston and slidable piston flange to thrust into the plunger shatter plate to shatter. The shattering of the plunger shatter plate causes the slidable piston and needle cannula to be thrust into the hollow plunger by the spring and is thus prevented from re-entering the guide tube. Again, no successful commercial product has resulted from this disclosure.\nU.S. Pat. No. 5,180,370 discloses a syringe which has an internal mechanism for retracting the needle into the syringe after the injection has been given. In one disclosed embodiment, the needle is manually retracted by pulling back on the plunger, and in another, the needle is propelled by a compressed spring into a hollow chamber within the plunger. A syringe produced with this disclosure would be complex to manufacture, and no successful commercial product has resulted from this disclosure.\nU.S. Pat. 5,188,599 discloses a hypodermic injection system with a needle that retracts within an interior cavity of the syringe plunger. The needle when retracted is held within the plunger. The disclosed device includes a cylindrical spring housing that has resilient fingers which capture a spring under bias holding a needle holder against the retaining force of resilient fingers. The plunger has a frangible end which dissociates when the outwardly tapered shoulders spread the resilient fingers, allowing the coiled spring to eject the needle and its holder into the interior cavity of the syringe plunger. The patent also discloses a body fluid sampling device that includes a double-ended needle for communication with an evacuated blood collection tube. This patent also includes a review of several earlier disclosures related to retractable needles. Attempts have been made to produce commercial products based on the disclosures of this patent, but as yet there is no successful commercial product.\nU.S. Pat. No. 5,201,710 discloses a syringe fitted with a clamping device for the needle and with a mechanism to enable the needle to be automatically retractable into the syringe body at the end of an injection. The disclosed device includes inner and outer cylinders, openings at the ends of the outer cylinder, a third opening at an end of the inner cylinder and a closure for the third opening. The disclosed device further includes a needle with a head, a seal, a first spring to push the needle against the closure and a clamping device loaded by a second spring to maintain outward to the syringe and to release the needle. There is a diaphragm in the closure that bends before breaking and a sharp element to break the diaphragm. There also is a closure to prevent the needle from being accessible and a stop to prevent the second cylinder from being moved outwardly after the syringe is used. As is apparent from the description, the device disclosed by this patent is complex and would be difficult to assemble. No successful commercial product has resulted from the disclosure in this patent.\nU.S. Pat. No. 5,385,551 discloses a non-reusable medical device that has a needle which is retractable by depression a plunger slidably mounted in the device. The disclosed device includes a front-mounted retraction mechanism that has a needle holder connected to the needle. The needle holder is supported along the axis of the device by a frictionally engaged retainer ring member coupled to the needle holder along an axially aligned sliding interface. The needle holder and retainer are positioned in the front portion of a hollow body. The front of a movable member or plunger presses against the retainer member passing around the needle holder which cannot move forward, thereby separating the retainer from the needle holder. The separation occurs by gradually reducing the extent of the sliding interface area until the retainer member pops loose from the needle holder whereupon the needle holder and needle are retracted into a cavity in the plunger in response to a retraction force applied to the needle holder by a previously compressed spring. Again, the device disclosed in this patent is complex, difficult to manufacture and appears to have significant undeliverable dead volume. Attempts have been made to commercialize products from this disclosure with only limited success.\nU.S. Pat. No. 5,407,436 discloses a hypodermic syringe that has a hollow needle that is automatically retractable after use. The disclosed syringe includes a one-piece body molding has a main chamber for a plunger, sample container or drug cartridge, a forward chamber to house a spring to bias a needle holder, and internal latching formations to retain the needle holder with the spring compressed in the forward chamber until automatic retraction when the latching formations are released by end of plunger movement. The patent discloses that the sealing between the plunger and the body is accomplished by an over-sized plunger head that forces head and wall deformation. The disclosed spring has seals at both ends for the forward chamber. The patent teaches that the needle, its holder, spring and seals can be installed using a sliding guide. In using a syringe produced using this disclosure, the practitioner would need to exercise care when drawing and expelling a fluid during filling, because the retraction of the needle is activated by depressing the plunger sufficiently to engage cooperating latches. The engagement occurs at the bottom of the stroke to expel fluid from the syringe.\nU.S. Pat. No. 5,769,822 discloses a non-reusable syringe with a hollow plunger that has a seal member thereon. The position of the plunger and the seal relative to the barrel permits the plunger, with sufficient strength, to carry applied pressure through the device during injection of a fluid and yet permit the seal disposed at one end of the plunger to have maximum sealing integrity between the plunger and a cylindrical barrel disposed around the exterior of the plunger to abate leakage of the liquid in a chamber within the barrel, as the plunger is manipulated from an expanded position to and expended position and thereafter to a third or collapsed position.\nU.S. Pat. No. 6,010,486 discloses a retracting needle syringe that substantially prevents reuse of the syringe by destroying the plunger rod and the needle hub and additionally, retracts the needle into the plunger rod. The disclosed syringe includes provisions that upon fully depressing the plunger rod and applying distally directed axial force, a frangible portion of the inner hub is broken and the plunger tip dislodges to allow a spring to urge a cutter to open the chamber inside the plunger.\nMost of the devices discussed in the above referenced disclosures are somewhat complex, and many require manufacture and assembly of parts with potentially difficult assembly or tight tolerance requirements. Many of the designs depend upon a careful application of forces by the practitioner to draw and expel fluids from the syringe. Also, if the tolerances between the multiple components of the device are not carefully adhered to during manufacture and assembly, normal usage may result in premature activation of the retraction function of the syringe. Current conventional syringes are considered by users to be virtually fault-free and reliable. They are used for a variety of different procedures involving both xe2x80x9cone-shotxe2x80x9d fill and inject procedures, as well as more complex mixing, measuring and delivery functions. In order for a retractable syringe to displace these functional, utilitarian and reliable conventional syringes, the retractable syringe should not significantly interfere with the users current practices, it needs to be substantially reliable and its cost should not be prohibitive. Current conventional syringes are often manufactured at rates of several hundred per minute and their cost is generally not a significant factor in their usage. Additionally every year, hundreds of millions of small capacity (one milliliter) syringes are used outside of the normal controlled health care environment by diabetics and other self-injectors who must daily accurately inject small amounts, often only a few tenths of a milliliter. These small capacity syringes are physically quite small, with an overall length of less than five inches and an inside bore diameter of less than one-quarter inch. Reviewing the disclosures above, one skilled in the art of high volume manufacturing recognizes that assembling hundreds of millions of most of these relatively complex devices with their retraction elements contained in such a small space as a one-quarter inch diameter bore is a daunting task. Additionally, many of the disclosed devices have substantial undeliverable xe2x80x9cdead volumesxe2x80x9d that substantially confound many diabetics\"\" need for accurate measuring, mixing of more than one type of insulin in the syringe and delivering small doses of insulin. The need thus exists for a selectively retractable syringe that is compatible with a small capacity syringe, that is capable of being manufactured at high volumes and is sufficiently non-complex to be reliable in use when produced at volumes of hundreds of millions per year. Such a device is disclosed herein below.\nA hypodermic syringe with a selectively retractable needle of the present invention includes an elongate barrel having an open proximal end and a distal end that defines a receiver with an inwardly projecting shoulder. The barrel has a hollow bore therethrough with an inside surface extending from the proximal end to the distal end. The syringe of the invention has a hollow elongate plunger with an open proximal end and a closed distal end. There is an elongate plug extending distally into the hollow plunger from the open proximal end forming an enclosed cavity within the plunger. The distal end of the plunger forms a slidable seal with the inside surface of the barrel to define a chamber for drawing and expelling fluid. The syringe also has an elongate hub having a passageway therethrough, a distally extending stem, and a proximal flange with an engagement for engaging the barrel. The stem is disposed within and sized for slidable movement within the receiver at the distal end of the barrel. The flange has a distal surface with a groove therein and a proximal surface defining the chamber in the barrel. The syringe of the invention has an elongate needle having a fluid path therethrough, a sharpened distal end and a proximal end. The needle is mounted in the passageway in the hub so that the sharpened distal end projects distally outwardly and the fluid path is fluidly communicative with the chamber in the barrel. The syringe also has an elongate spring disposed about the stem of the hub that is compressed between the flange and the inwardly projecting shoulder of the receiver to provide a bias. There is a hollow sleeve sized to fit within the receiver over the elongate spring. The sleeve has a distal end disposed at shoulder and a proximal end with a sharpened edge that is disposed in the groove in the distal surface of the flange. When a user applies a sufficient force, a force greater than necessary to expel fluid from the chamber, to the plunger, the hub is sufficiently moved distally in the receiver for the cutting edge of the sleeve to cut through the flange and the closed distal end of the plunger to expose the cavity in the plunger. When the cavity in the plunger is exposed, the bias of the spring urges a sufficient movement of a cut portion of the hub having the needle mounted therein, a cut portion of the distal end of the plunger and the sleeve into the cavity in the plunger to a position wherein an inadvertent exposure of the sharpened distal point is substantially prevented.\nThe syringe of the invention has an undeliverable xe2x80x9cdead-spacexe2x80x9d volume substantially similar to fixed needle xe2x80x9clow dead-spacexe2x80x9d conventional syringes, i.e., substantially no undeliverable volume. The syringe of the invention is as suitable for use in drawing, measuring, mixing and delivering small volumes of medicaments as conventional syringes. Unlike many of the devices disclosed above, the syringe of the invention is substantially unlikely to be inadvertently retracted by a user following currently used practices and procedures. The syringe of the invention does not depend on a user having to exercise substantially more care than with a conventional syringe when drawing and mixing fluids in the syringe to avoid inadvertent activation, and importantly, the syringe of the invention is compatible with the efficiency of high volume automated manufacture that utilizes much existing manufacturing equipment. Once needle is retracted in the syringe of the invention, the syringe cannot be restored to functionality, as the hub flange is cut through and the plunger is cut through rendering the syringe substantially unusable and protecting the needle point from inadvertent contact by anyone.\nAn alternate embodiment of the syringe of the present invention includes an elongate barrel having an open proximal end and a distal end defining a receiver. The barrel also includes a hollow bore therethrough. A hollow elongate plunger having a proximal end and a distal end defining a cavity therein is provided. The distal end of the plunger forms a slidable seal with the inside surface of the barrel. An elongate hub comprises a passageway therethrough including a distally extending stem, and a proximal flange with an engagement for engaging the barrel. The stem being disposed within and sized for movement in the receiver at the distal end of the barrel. The flange of the hub includes a distal surface and a proximal surface. An elongate needle having a fluid path therethrough is connected to the stem of the hub so that the passageway is in fluid communication with the chamber of the barrel. An elongate spring is disposed about the stem of the hub. The spring is compressed to provide a biasing force to the flange. A hollow sleeve sized to fit within the receiver includes a distal end and a proximal end having a cutting edge facing the distal surface of the flange of the hub so that when the user supplies a sufficient force, greater than a force necessary to expel fluid from the chamber, to the plunger, the hub is sufficiently moved distally in the receiver for the cutting edge of the sleeve to cut through the flange and the closed distal end of the plunger to expose the cavity therein. This allows the bias of the spring to urge a sufficient proximal movement of the cut portion of the hub having the needle mounted therein and the cut portion of the distal end of the plunger into the cavity. In the event that the frictional force between the inside of the sleeve and the cut portions of the hub and the plunger is greater than the frictional force between outside of the sleeve and the remaining portions of the hub and the plunger, the spring will move the sleeve from its position in the receiver so that the cut portion of the hub, the cut portion of the distal end of the plunger and the sleeve move into the cavity. If the frictional force on the outside of the sleeve is greater than the frictional force on the inside of the sleeve the sleeve will remain in its position in the receiver while the cut portion of the hub and the cut portion of the distal end of the plunger move into the cavity."}
{"text": "It is well accepted that submersible automated pool cleaning devices, such as self propelled pool cleaning vehicles (PCV's) are essential to the proper maintenance of a pool, whether the pool be above or below ground. The typical vehicle includes a housing and drive members. The drive members attach to the housing usually through connection to a chassis. Drive members include wheels, endless loop tracks and combinations of each. Additionally, vehicles have included two and four wheel drive vehicles in various combinations and variations. In the case of a belt or endless loop track, the track wraps around the drive and/or idler wheels or rollers.\nIt will be appreciated that such vehicles may also have some combination of wheels and/or rollers and the vehicle may also be a two or four-wheel drive vehicle. The endless loop track may be fitted over any combination of two or four wheel drive vehicles. Such tracks have been found to be an effective means for moving the vehicle around the surface of the pool while the vehicle is submerged.\nOf course, the primary purpose for all PCV's is to clean the pool and pool water. In order to properly maintain clean, pool water, the water itself as well as the pool surfaces must be cleaned and kept clean. It is imperative to keep the pool surfaces free from the build up of dirt and debris. Once dirt and/or debris is allowed to settle in on the pool surface, algae forms. Marge enough build-up of such algae can cause the pool pH to become unbalanced and the pool water unstable, in terms of cleanliness.\nSometimes drastic efforts caused by “dirty” water result in the entire contents of the pool being drained and starting again. As can be appreciated this is an expensive and time consuming process. Additionally, without thorough cleaning and maintain, the situation will simply be repeated. There is no guarantee, without proper maintenance methods, that the water will not simply revert to the same “dirty” condition.\nOf special concern in this regard, is the fact that embedded particles can be particular concern. As noted by others, these embedded particulates can cause fatal consequences in term of “dirty water” requiring the emptying of the pool for dry cleaning. Clearly, a condition that is primary to avoid.\nClearly, there is a strong industry-wide need to prevent the build up of particulates in the pool water. There is even a stronger need to prevent the embedding of such particulates in the surface of the pool. Cleaning each and every square millimeter of the pool's surface is simply not a practical solution. Therefore, Applicant has developed a structure for using the automated pool cleaning vehicle to remove even embedded particulates in the surface of the pool. Additionally, the cleaning structure, in accordance with the invention, facilitates the PCV to thoroughly clean the dirt and debris thus dislodged from the pool's surface. This not only serves the short term goal of keeping the water clean and free from dirt and debris, but also is helpful in accomplishing the long term goal of preventing the build up of such particulates on and in the surface of the pool.\nWhat is needed is cleaning structure for a submersible automated pool cleaning vehicle (PCV), which is adapted to clean the pool surface by dislodging dirt and debris from the surface of the pool and further to dislodge even embedded dirt and debris from the surface of the pool using the automated motion of the PCV."}
{"text": "Schizophrenia is a progressive and devastating neurological disease characterized by episodic positive symptoms such as delusions, hallucinations, thought disorders and psychosis and persistent negative symptoms such as flattened affect, impaired attention and social withdrawal, and cognitive impairments (Lewis D A and Lieberman J A, Neuron, 2000, 28:325-33). For decades research has focused on the “dopaminergic hyperactivity” hypothesis which has led to therapeutic interventions involving blockade of the dopaminergic system (Vandenberg R J and Aubrey K R., Exp. Opin. Ther. Targets, 2001, 5(4): 507-518; Nakazato A and Okuyama S, et al., 2000, Exp. Opin. Ther. Patents, 10(1): 75-98). This pharmacological approach poorly address negative and cognitive symptoms which are the best predictors of functional outcome (Sharma T., Br. J. Psychiatry, 1999, 174(suppl. 28): 44-51).\nA complementary model of schizophrenia was proposed in the mid-1960′ based upon the psychotomimetic action caused by the blockade of the glutamate system by compounds like phencyclidine (PCP) and related agents (ketamine) which are non-competitive NMDA receptor antagonists. Interestingly in healthy volunteers, PCP-induced psychotomimetic action incorporates positive and negative symptoms as well as cognitive dysfunction, thus closely resembling schizophrenia in patients (Javitt D C et al., 1999, Biol. Psychiatry, 45: 668-679 and refs. herein). Furthermore transgenic mice expressing reduced levels of the NMDAR1 subunit displays behavioral abnormalities similar to those observed in pharmacologically induced models of schizophrenia, supporting a model in which reduced NMDA receptor activity results in schizophrenia-like behavior (Mohn A R et al., 1999, Cell, 98: 427-236).\nGlutamate neurotransmission, in particular NMDA receptor activity, plays a critical role in synaptic plasticity, learning and memory, such as the NMDA receptors appears to serve as a graded switch for gating the threshold of synaptic plasticity and memory formation (Hebb D O, 1949, The organization of behavior, Wiley, NY; Bliss TV and Collingridge G L, 1993, Nature, 361: 31-39). Transgenic mice overexpressing the NMDA NR2B subunit exhibit enhanced synaptic plasticity and superior ability in learning and memory (Tang J P et al., 1999, Nature: 401-63-69).\nThus, if a glutamate deficit is implicate in the pathophysiology of schizophrenia, enhancing glutamate transmission, in particular via NMDA receptor activation, would be predicted to produce both anti-psychotic and cognitive enhancing effects.\nThe amino acid glycine is known to have at least two important functions in the CNS. It acts as an inhibitory amino acid, binding to strychnine sensitive glycine receptors, and it also influences excitatory activity, acting as an essential co-agonist with glutamate for N-methyl-D-aspartate (NMDA) receptor function. While glutamate is released in an activity-dependent manner from synaptic terminals, glycine is apparently present at a more constant level and seems to modulate/control the receptor for its response to glutamate.\nOne of the most effective ways to control synaptic concentrations of neurotransmitter is to influence their re-uptake at the synapses. Neurotransmitter transporters by removing neurotransmitters from the extracellular space, can control their extracellular lifetime and thereby modulate the magnitude of the synaptic transmission (Gainetdinov R R et al, 2002, Trends in Pharm. Sci., 23(8): 367-373).\nGlycine transporters, which form part of the sodium and chloride family of neurotransmitter transporters, play an important role in the termination of post-synaptic glycinergic actions and maintenance of low extracellular glycine concentration by re-uptake of glycine into presynaptic nerve terminals and surrounding fine glial processes.\nTwo distinct glycine transporter genes have been cloned (GlyT-1 and GlyT-2) from mammalian brain, which give rise to two transporters with ˜50% amino acid sequence homology. GlyT-1 presents four isoforms arising from alternative splicing and alternative promoter usage (1a, 1b, 1c and 1d). Only two of these isoforms have been found in rodent brain (GlyT-1a and GlyT-1b). GlyT-2 also presents some degree of heterogeneity. Two GlyT-2 isoforms (2a and 2b) have been identified in rodent brains. GlyT-1 is known to be located in CNS and in peripheral tissues, whereas GlyT-2 is specific to the CNS. GlyT-1 has a predominantly glial distribution and is found not only in areas corresponding to strychnine sensitive glycine receptors but also outside these areas, where it has been postulated to be involved in modulation of NMDA receptor function (Lopez-Corcuera B et al., 2001, Mol. Mem. Biol., 18: 13-20). Thus, one strategy to enhance NMDA receptor activity is to elevate the glycine concentration in the local microenvironment of synaptic NMDA receptors by inhibition of GlyT-1 transporter (Bergereon R. Et al., 1998, Proc. Natl. Acad. Sci. USA, 95: 15730-15734; Chen L et al., 2003, J. Neurophysiol., 89 (2): 691-703).\nGlycine transporters inhibitors are suitable for the treatment of neuroligical and neuropsychiatric disorders. The majority of diseases states implicated are psychoses, schizophrenia (Armer R E and Miller D J, 2001, Exp. Opin. Ther. Patents, 11 (4): 563-572), psychotic mood disorders such as severe major depressive disorder, mood disorders associated with psychotic disorders such as acute mania or depression associated with bipolar disorders and mood disorders associated with schizophrenia, (Pralong E T et al., 2002, Prog. Neurobiol., 67: 173-202), autistic disorders (Carlsson M L, 1998, J. Neural Transm. 105: 525-535), cognitive disorders such as dementias, including age related dementia and senile dementia of the Alzheimer type, memory disorders in a mammal, including a human, attention deficit disorders and pain (Armer R E and Miller D J, 2001, Exp. Opin. Ther. Patents, 11 (4): 563-572).\nThus, increasing activation of NMDA receptors via GlyT-1 inhibition may lead to agents that treat psychosis, schizophrenia, dementia and other diseases in which cognitive processes are impaired, such as attention deficit disorders or Alzheimer's disease."}
{"text": "Typical natural gas fired engines, such as those used in the petroleum industry to compress natural gas at well heads and along pipelines, have up to 16 cylinders, often generating between 500 to 3000 HP. These engines normally are run continuously near full load conditions with shut downs primarily being for maintenance such as for oil changes. The continuous operation near full load, of course, places severe demands on the engine lubricant. Indeed, because the lubricant is subjected to a high temperature environment oxidation processes occur which limit lubricant life. Also, natural gas engines emit nitrogen oxides (NOx) some of which comes into contact with the lubricant resulting in nitration processes that also will limit lubricant life. Typically these processes are accompanied by increases in oil viscosity. Thus it is desirable to extend the life of gas engine oils by enhancing the oil's resistance to oxidation and nitration and to reduce viscosity increases in the oil.\nTo extend lubricant life, base oils are formulated with various additives such as dispersants, detergents, antioxidants, viscosity index improvers and the like, to provide a lubricating oil composition. This art of lubricating oil formulation, however, has become increasingly complex with ever more stringent requirements of end-users. Indeed, experience has shown that incorporation of one type of additive in a lubricant composition can have a negative impact on the function of another type of additive. Consequently extensive research continues in the quest for lubricants of improved life and function. Additionally, gas engine oil lubricants are typically formulated with metal containing detergents and metal containing antiwear additives to provide lubricant ash that functions to protect the valve and seat interface of the engine. The ability to provide an ashless gas engine oil, i.e., one having an ash level as determined by ASTM test method D874 below 0.1 mass %, that has enhanced life is particularly difficult.\nOne object of the present invention is to provide an ashless gas engine lubricating composition that has enhanced resistance to oxidation and nitration.\nAnother object of the invention is to provide an ashless gas engine oil lubricating composition that has improved life as evidenced by reduction in viscosity increase."}
{"text": "Wireless communication networks are widely deployed to provide various communication services such as telephony, video, data, messaging, broadcasts, and so on. Such networks, which are usually multiple access networks, support communications for multiple users by sharing the available network resources. One example of such a network is the UMTS Terrestrial Radio Access Network (UTRAN). The UTRAN is the radio access network (RAN) defined as a part of the Universal Mobile Telecommunications System (UMTS), a third generation (3G) mobile phone technology supported by the 3rd Generation Partnership Project (3GPP). The UNITS, which is the successor to Global System for Mobile Communications (GSM) technologies, currently supports various air interface standards, such as Wideband-Code Division Multiple Access (W-CDMA), Time Division-Code Division Multiple Access (TD-CDMA), and Time Division-Synchronous Code Division Multiple Access (TD-SCDMA). The UMTS also supports enhanced 3G data communications protocols, such as High Speed Packet Access (HSPA), which provides higher data transfer speeds and capacity to associated UMTS networks.\nMulti-subscriber identity module (SIM) wireless devices have become increasing popular because of their flexibility in service options and other features. One type of multi-SIM wireless device, a dual-SIM dual-standby (DSDS) device, enables both SIMs to be in idle mode waiting to begin communications, but only allows one SIM at a time to participate in an active communication due to sharing of a single radio frequency (RF) resource (e.g., transceiver). Other multi-SIM devices that extend this capability to more than two SIMs are referred to as multi-SIM, multi-standby (MSMS) devices, and may be configured with any number of SIMs greater than two (i.e., multi-SIM multi-standby wireless devices).\nMulti-SIM wireless devices can have multiple subscriptions to one or more wireless networks. For example, in a DSDS device, a first subscription may support networks using a first technology standard, such as Long Term Evolution (LTE) or Wideband Code Division Multiple Access (WCDMA), while a second subscription may support a second technology standard, such as Global System for Mobile Communications (GSM) Enhanced Data rates for GSM Evolution (EDGE) (also referred to as GERAN). A DSDS device enables both SIMs to be in idle mode waiting to begin communications on their respective networks, but only allows one SIM at a time to participate in an active communication due to sharing of a single RF resource (e.g., a transceiver) configured to both communication technologies."}
{"text": "Three-dimensional (3-D) object printers, such as those which employ Fusion Deposition Modeling (FDM), are known. The printing process for such a device involves the deposition of printing material onto a printing platform, also referred to as a print bed. The deposited material may be melted into a pliable state, extruded through a heated nozzle and built up, layer by layer, until the final result is a 3-D object. Because the layers are deposited in sequence on top of each other, print success and quality depend upon the ability to maintain registration of the object with the extruder nozzle throughout the duration of the print job to ensure that each stacked layer registers with the previous one.\nPrint success and quality may also depend upon adequate adhesion between the printed object and the print bed. Sometimes the first few layers of the printed object do not have sufficient adherence to the print bed, causing the printed object to shift, warp, or delaminate from the print bed, resulting in a failed or poor quality printed object. The print beds for known FDM style 3-D printers are typically made of metal, glass or acrylic. These print beds are not considered consumables, nor are they ideally suited to provide reliable surfaces on which the 3-D printed objects can adhere solidly and consistently. Therefore, it is preferable to pretreat and/or cover the print bed surface of a FDM style 3-D printer prior to printing an object so as to prevent damaging the print bed and to improve the likelihood that the printed object will adhere adequately to the print bed for the duration of the print.\nManufacturers and users of FDM style 3-D printers often require or recommend covering the print bed surface with heat resistant polyimide film or paper masking tape (typically used by painters), and/or pre-treating the surface with hairspray, special water-soluble glue or other liquid treatment solutions. These applications are intended to hold the 3-D printed object to the surface while it is printing and to preserve the longevity of the original print bed. These print bed covering materials are intended to provide a removable and replaceable surface on which to print, and in some cases take the wear and tear that would otherwise be inflicted upon the print bed.\nPrinter manufacturers may also recommend using a heated print bed for 3-D printing deposit material that requires slower cooling time, such as ABS (acrylonitrile butadiene styrene). Often times the heated print bed is heated to temperature up to 100-130 degrees Celsius. In such cases, the print bed cover material must be able to function while be exposed to the elevated temperature of the print bed. Accordingly, heat resistant polyimide films may be used for heated print beds, while paper masking tape cannot. Other types of deposit material, such as PLA (polylactide), are not slow cooling and do not require a heat tolerant covering material, such as paper masking tape. In general these two common alternatives (heat resistant polyimide films and paper masking tapes) cannot be used interchangeably with different deposit materials due to their differing properties and heat resistant limitations. Thus, there is a need for a single type of print bed cover which may be used with both heated and non-heated print beds.\nFurther, there are number of disadvantages that may arise from using known print bed covers. For example, commonly used heat resistant polyimide films or paper masking tapes may be difficult and tedious to apply or install on the print bed. Polyimide films and paper masking tape generally are supplied on a roll and need to be cut and resized for the print bed on which they are installed. If the width of the supplied roll is not as wide as the print bed, then multiple sheets of the film or tape may need to be applied side by side in order to cover the print bed. However, it is extremely important that the print bed surface be flat and level to the extruder nozzle; i.e., the gap between the extruder nozzle and top surface of the print bed needs to be uniform over the entirety of the print bed. Failure to provide a uniform distance between the extruder and the print bed cover may result in defective print objects and even damage or tearing of the cover if the extruder nozzle contacts it. Therefore, for best 3-D printing results, the films and tapes need to be applied without overlapping seams, folds, creases or air bubbles under or in the covering surface since such irregularities may cause variation of the distance between the extruder nozzle and the print surface. Because the known films and tapes are typically very thin, they are difficult to work with in a manner that avoids overlapping and air bubbles, and are susceptible to unwanted stretching, folding and creasing while being adhered to the print bed.\nAnother difficulty which arises from using polyimide film or paper masking tape is that they may not provide sufficient adhesion to keep the 3-D printed object from moving or warping upward during the printing process, resulting in a failed or defective printed object. Alternatively, in some instances, polyimide film or paper masking tape provide so much adhesion that upon completion, the 3-D printed object is difficult to remove from the print surface, which can result in damage to the print bed cover or printing device, or even in personal injury. Excessive adhesion may be further complicated by the thinness of the polyimide film and paper masking tape. When adhesion is too great, the film or tape may be damaged when the printed object is removed from it, or when leftover deposited material is scraped off. The foregoing challenges may result in the need for frequent replacement of the polyimide film or paper masking tape.\nWater-soluble glue is another known covering material for 3-D printer beds. The use of such glues, however, requires exposing the completed 3-D printed object to water while it is still adhered to the print surface (through various methods) in order to separate the 3-D printed object from the print surface. This water application step can be inefficient since (a) it is an extra step that would not otherwise be required, (b) requires a new application of glue for each printing operation, and (c) may not work in many situations where the original print bed is equipped with an FDM style 3-D printer bed that is not easily removable from the printing device. Extremely rigid perfboard or FR4 fiberglass (which is typically used to make printed circuit boards) are still other covering materials that are currently used as print surfaces or covers. Such materials may also suffer from removal problems."}
{"text": "The expandable sealants, which are used in automobiles to fill gaps and voids as well as improve the acoustical properties, have certain inherent problems. These materials at present must be applied at ambient temperatures so as not to prematurely cure or expand the sealant. In this condition, the sealant does not easily flow and is therefore contoured to a specific configuration prior to its application to the automobile. This results in certain instances, in a requirement that as many as twenty-two different sealant configurations must be formed and inventoried for one particular vehicle.\nAnother drawback to these \"cold applied\" sealants is that they do not adhere sufficiently to the vehicle and periodically are knocked off or moved from their original location by subsequent cleansing and priming tanks through which the vehicle must pass prior to the sealant being expanded and cured thereby securely positioning it in place. This results in contamination of the cleansing tanks and processes as well as an inferior automobile.\nOne attempt to overcome these drawbacks might be to apply the sealant hot. Such a hot applied pumpable sealant must not begin to cross link or blow at the temperatures at which the material is pumped as this will create problems in the pump as well as difficulty controlling the quality of the final sealant.\nTherefore, what is needed in the art is in expandable sealant which is easily applied, will not be affected by the cleaning or painting processes, and does not activate prematurely at application temperatures."}
{"text": "Polished fingernails are a necessity for many people. In addition, people ranging in age from their early teens up to when they become elderly engage in such grooming activity. Professional manicures also can be expensive for some people, while others merely may wish to paint and polish their own fingernails. The present invention is geared to those who want to care for their own fingernails in an organized, comfortable, clean and efficient manner.\nTypically, many people possess a bag or carrying case that contains the multitude of fingernail polish. This means that many bottles featuring varying colors and brands are stuffed into at least one bag. In addition, other related items dealing with such aspects as cuticles, filing and clipping also are scattered throughout. This creates a problem of organization, as the bag often must be emptied of its contents so that a particular color or brand can be utilized. The present invention solves this problem by forming holes into an organizational surface that are elevated and supported by intermittently placed platforms. The bottles can be placed into the holes. Meanwhile, because the platforms elevate the organizational surface, the placement of these platforms permits the user a clear line of sight to instantly identify the varying colors and brands of the bottle instead of just the bottle cap.\nPainting and polishing fingernails also requires a degree of accuracy and stability. For example, users typically must paint their fingernails while holding their hand in the air. Some people find that their hands shake and require the hand to be rested on a table or other fixed surface. Others merely find that they need a surface to keep their hand from moving as they concentrate on the act of painting or polishing. The present invention solves this problem by including an elevated ramp that extends from the base of the present invention. The preferred embodiment of the elevated ramp extends at a 45-degree angle. The elevated ramp allows the user to rest his or her hand at an angle so that the grooming process can be achieved in an accurate and stable manner. The present invention also accommodates users facing the awkwardness of comfortably grooming their fingernails by forming a curvature into the front of the base of the present invention. The curvature is such that the present invention is more stable and close to the user's body as it wraps partially around the torso or lap.\nU.S. Pat. No. 5,913,312 issued to Donnell on Jun. 22, 1999 is a nail polish bottle holder. Unlike the present invention, Donnell is formed from a sheet of flexible foam material with designated areas available so that polish bottles can be pushed into the foam. The present invention, in contrast, provides holes in such a manner that fingernail-grooming bottles can seamlessly be placed into a hole so that the bottle label or color is still easily identifiable. In addition, Donnell requires that a user's hand be rested on the surface of the foam. However, the present invention provides an elevated ramp for this purpose. The elevated ramp, along with other features of the present invention, can be securely fitted onto a base of the present invention. Meanwhile, a lip for spillage protection extends along the front of the base, the front of the base curved in such a manner as to wrap around the front of a person's body or around a lap. Donnell, in contrast, does not account for shape or angle that add to such important areas as stability, containment and ease of use."}
{"text": "5th generation (5G) mobile communications, initial commercialization of which is expected around 2020, is recently gathering increased momentum with all the worldwide technical activities on the various candidate technologies from industry and academia. The candidate enablers for the 5G mobile communications include massive antenna technologies, from legacy cellular frequency bands up to high frequencies, to provide beamforming gain and support increased capacity, new waveform (e.g., a new radio access technology (RAT)) to flexibly accommodate various services/applications with different requirements, new multiple access schemes to support massive connections, and so on. The International Telecommunication Union (ITU) has categorized the usage scenarios for international mobile telecommunications (IMT) for 2020 and beyond into 3 main groups such as enhanced mobile broadband, massive machine type communications (MTC), and ultra-reliable and low latency communications. While all the requirements need not be met simultaneously, the design of 5G networks may provide flexibility to support various applications meeting part of the above requirements on a use case basis."}
{"text": "1. Field of the Invention\nThe invention relates to a transmission for a motor vehicle, and more particularly to a transmission comprising an input shaft, a countershaft, and an output shaft arranged concentrically with the input shaft and the output shaft journalled to the input shaft, a pair of gearwheels arranged between the input or output shaft, and the countershaft constituting a constant ratio, and a plurality of pairs of gearwheels carried on the gear shafts and each comprising a fixed gearwheel and a loose gearwheel for engaging the various gear steps.\n2. Description of the Prior Art\nA prior art transmission is shown in German specification 34 45 519 (DE '519). As shown in FIG. 1 of the DE '519 patent, the input shaft and output shaft and/or the countershaft are mounted in three bearings and a synchronizer unit for engaging the 5th gear and the reverse gear is provided on the countershaft.\nThe transmission in the DE '519 patent has the disadvantage that the gearwheel pair for the constant ratio is arranged directly on a short input shaft and at the input end of the countershaft for all the following gear steps. This results in the disadvantage that when the transmission is in the neutral position with the friction clutch engaged, the input shaft constantly drives the countershaft. Thus, by way of the fixed gearwheels carried on the countershaft, the corresponding loose gearwheels on the output shaft are also driven. This leads to undesirable rattling noises in neutral, which arise from the various imperfections in the gear teeth and play between the teeth.\nGerman specification 37 29 818 (DE '818) discloses a similar transmission, wherein a synchronizer unit for engaging the 5th gear and the synchronized reverse gear is carried on the output shaft.\nThe transmission of the DE '818 patent again has the disadvantage that the constant ratio is arranged directly between a short input shaft and the input end of the countershaft, and accordingly in the neutral position of the transmission with the friction clutch engaged the whole set of gearwheels on the countershaft rotates, whereby, owing to the loose gearwheels being in mesh with the fixed gearwheels, and therefore marked rattling noises can be caused.\nIt would therefore be desirable to provide a transmission having a neutral position wherein the gear ratios are not driven and therefore less gear noise occurs. A further object is to provide six forward gears and a synchronized reverse in such a transmission."}
{"text": "1. Field of the Invention\nThe present invention relates to a data providing system and a computer-readable recording medium storing data providing program for providing electronic data. In addition, the present invention relates to a data receiving system and a computer-readable recording medium storing data receiving program for receiving electronic data.\n2. Description of the Related Art\nIn recent years, when a data provider offers electronic data such as a computer program, the following method has been generally used. For example, a program is stored in a connectable server on the Internet and downloaded to a user (i.e., a data receiver) who connects with this. Besides, an electronic device possessed by a user as the data receiver is not limited to dedicated equipment having only a specific function (e.g., a copying function). As the market has been increasingly satisfied with hardware resources, an electronic device which has diverse functions complexly has been popular. Therefore, even after an electronic device has been newly introduced, if additional functions are incorporated therein, the electronic device can be given a multiple function and customized. This is serviceable to a user.\nThe above described background has recently encouraged the art of equipping an image forming apparatus such as a copying machine with an execution environment for a Web application program and installing an external application program in an image forming apparatus on the Internet. This prior art enables a user to access a Web page using a client terminal (i.e., a PC) linked to the Internet, and via this Web page, install the external application program in the image forming apparatus. This external application program is stored, for example, into a server, an IC card or the like which is linked to the Internet.\nHowever, according to such an aspect of the prior art in which the external application program is stored in the server and installed from there in the image forming apparatus, a disadvantage arises in that a user tends to suffer damage, such as illegitimate read-out of data stored in the server and tapping of communication with the server. There is also a disadvantage in that the contents of the external application program may be maliciously altered through unjust access to the server and the application program altered without being noticed can be installed in the image forming apparatus.\nIn addition, if a third party with evil intent establishes a server and pretends to be a distributor of an application program, a user may install the application program from there unwittingly."}
{"text": "The biochemical processes that occur during blood preservation or storage contribute to the diminution of the post-transfusion viability, i.e. survival of stored rec cells after transfusion to a recipient, and oxygen off-loading capacity of red cells, both of which are statistically correlated with the duration of storage period. Previous studies of stored red blood cells have indicated that intracellular levels of ATP (adenosine triphosphate) and 2,3-DPG (2,3-diphosphoglycerate) largely or entirely determine the post-transfusion viability and oxygen off-loading capabilities respectively of the red cell. It is well-established that both of these substances are products of glycolysis, a biochemical pathway that degrades glucose.\nMaintenance of post-transfusion viability of stored red cells is closely correlated with the levels of cellular ATP, a high-energy compound. Glycolysis is the only ATP generated pathway in the red cell. It appears that ATP preserves membrane integrity by maintaining proper ionic gradients across the red cell membrane, adequate lipid turnover rate, hemoglobin in a functional state and normal equilibrium of oxidized and reduced glutathione, along with synthesis of adequate amounts of NAD+and NADP+, i.e. Nicotinamide adenine dinucleotide and its phosphate.\nIt is well-established that oxygen off-loading ability of red cells is determined by another glycolytic intermediate, 2,3-DPG. This compound declines as a function of time during storage, most likely secondary to both a decreased rate of synthesis and an increased rate of degradation. The red cells with low 2,3-DPG show increased oxygen affinity or decreased ability to release oxygen at the tissue level. Thus, the stored red cells are less efficient vehicles of oxygen transport, the most important function of the red cells--stated otherwise, this means that stored and therefore 2,3-DPG depleted red cells are of poor quality with regard to their function.\nIt follows that it would be highly desirable to be able to maintain near-normal red cell 2,3-DPG and ATP as sucn a result would have profound effects in terms of both red cell function and post-transfusion survivability in vivo. Thus, adequate maintenance of both of these compounds would permit the storage of red cells for a longer period of time, which would alleviate problems of shortage of blood for transfusion and low quality of stored blood.\nPrevious studies of stored whole blood and stored red cell concentrates have indicated that intracellular levels of ATP and 2,3-DPG are important in extending storage capabilities. To this end, several studies have been conducted which have incorporated various chemical additives along with CPD (citrate-phosphate-dextrose) anticoagulant to stimulate glycolysis, yielding a net increase in red cell ATP level. One of the commercial additives that has recently been studied is adenine. The incorporation of adenine along with CPD anticoagulant into stored blood appears to increase ADP (adenine diphosphate) levels, thereby driving the glycolytic equilibrium toward the synthesis of ATP. However, adenine has an adverse effect on the maintenance of the level of 2,3-DPG, i.e. it lowers 2,3-DPG level with concomitant poor function of the stored red cells.\nRecent concern over the levels of ATP and 2,3-DPG has become a controversial subject. Because the main objective of transfusing patients is to provide or improve the oxygen delivery to the tissues, the blood oxygen affinity, directly determined by 2,3-DPG, is of critical importance. Therefore, in providing patents with suitable blood for transfusion, one must now consider not only red cell viability in vivo but also hemoglobin oxygen affinity for adequate oxygen transport function, the ultimate goal of red cell transfusion. As a result, research has also been geared towards incorporation of chemicals into the CPD and other preservative media to increase 2,3-DPG and ATP levels.\nThe significance of near-normal 2,3-DPG-containing red cells becomes self-evident when one examines various clinical conditions such as congestive heart failure, right to left cardiac shunts, and hypoxemia due to pulmonary disease, where patients singularly require the oxygen transport function of the transfused red cells. The transfused red cell, totally depleted of 2,3-DPG, is said to regain half the normal level of this substance within about 24 hours, but this increase may not be rapid enough to be effective in severely ill patients. Furthermore, it is not known whether the rate of resynthesis of 2,3-DPG in the donor cells given to a critically ill patient is comparable to that observed in normal recipients. Dennis et al. (Surgery 77 (6):741-747, June 1975) has reported a direct correlation between the ability to compensate for low 2,3-DPG levels and the severity of the illness of the patient. Blood with nearly normal hemoglobin-oxygen affinity is thus preferable for use in massive transfusions, particularly in infants, older patients, and patients with complicating cardiovascular and pulmonary disease.\nThe physiological effects of high oxygen-affinity (2,3-DPG depleted) red cells on the myocardial, cerebral, hepatic, and renal functions have not yet been fully evaluated, although patients requiring massive transfusions seem to be most susceptible to the adverse effects due to very low levels of 2,3-DPG; see Beutler et al., Vox Sang. 20:403-413 (1970).\nAlthough numerous investigations indicate that the levels of ATP and 2,3-DPG can be better maintained when the two chief preservative solutions ACD (acid citrate dextrose) and CPD (citrate phosphate dextrose) are supplemented with adenine, inosine, or both during storage at 4.degree. C., this conclusion must be approached with some caution. As has been reported by Bunn et al. in New England Journal of Medicine 282:1414-1421 (1970), a patient receiving three or four units of thus-supplemented blood may develop hyperuricemia, which persists for approximately 24 hours. As reported by Valeri in J. Med. (Basel) 5(5):278-291 (1974), a further cause for concern is the possible renal toxicity of 2,8-dioxyadenine, a metabolite of adenine. Moreover, additives that maintain ATP level, i.e. adenine, tend to lower 2,3-DPG level and those that maintain 2,3-DPG tend to lower ATP level, thus making the maintenance of both of these compounds a currently unrealizable goal. No matter which chemical is used with an ACD or CPD preservative solution, it appears that only a combination of various chemical additives will maintain 2,3DPG levels in blood under refrigeration conditions for an acceptable period of time.\nDieindoerfer et al, in U.S. Pat. No. 3,795,581, disclose a method of storing and preserving whole blood using an aqueous solution of dihydroxyacetone to increase the 2,3-DPG content. In a later patent, 3,874,384, Deindoerfer et al disclose the use of a combination of dihydroxyacetone and ascorbic acid to maintain DPG levels in stored blood. Estep, No. 4,386,069, discloses the use of a fatty ester having at least two ester linkages comprising fatty hydrocarbon groups of from four to twelve carbon atoms each to enhance the preservation of normal red blood cell morphology during storage. Harmening, No. 4,112,070, discloses a process for extending the useful shelf life of red cells by maintaining adequate levels of ATP and 2,3-DPG by adding an insoluble polymer material as a source of inorganic phosphate ions during the storage period. Harmening-Pettiglion, No. 4,390,619, discloses a method for extending the shelf life of blood platelets by maintaining both the pH and ATP levels suitable for transfusion. This is accomplished by providing to the platelets a water-insoluble polymer containing releasable weakly basic buffer ions capable of continuously supplying buffer to the platelets."}
{"text": "The present invention relates to a method for producing a glass ceramic having a garnet phase and to an advantageous application of such glass ceramic.\nRecently, LEDs have been employed to an increasing degree for lighting purposes because they offer a number of advantageous properties including, for example, high efficiency due to direct conversion of electric energy to light energy, and high compactness.\nHowever, until a few years ago LEDs were employed in “low-emitting” applications only, especially for indication purposes. The high potential of LEDs for applications with high light demand was discovered only recently when increased efforts were made to achieve improved energy input coupling and improved heat management.\nLEDs produce light in a very narrow spectral band, while in most of the cases white light is required for illumination purposes. Commercially available white LEDs use a III-V semiconductor emitter for stimulating a luminescent material that emits a secondary wavelength in a lower wavelength band (down conversion). One known solution uses a blue InGaN/GaN LED for stimulating YAG:Ce, a broadband yellow luminescent material. With these LEDs, that have been converted using a luminescent material, a given proportion of the emitted blue light passes the luminescent layer covering the LED chip so that the overall spectrum obtained assumes a color very close to white light. Due to the absence of spectral portions in the blue/green band and in the red wavelength band, the resulting color is, however, not satisfactory in most of the cases.\nUS 2003/0025449 A1 discloses a hybrid LED comprising a radiation-emitting semiconductor body (chip) which is in direct contact, via an optical coupling agent, with a glass ceramic body that serves as a conversion agent. The glass ceramic body contains crystallites of the garnet type doped with rare earths (such as YAG:CE), thiogalate or chlorosilicate as luminescent material. The starting glasses from which such glass ceramics are produced consist of silicate glasses or borate glasses. The luminescent glass ceramic is produced by mixing a glass frit with a suitable proportion of a luminescent material in powder form, and the mixture is molten, then cast and molded to achieve the desired shape. It is possible in this way, from the very beginning, to produce a glass ceramic body in the desired shape, advantageous for the intended application, for example in the form of a lens.\nHowever, that document does not disclose the way in which to produce such a luminescent glass ceramic having a garnet phase, with properties as advantageous as possible. Rather, the document merely relates to the melting technology used for production in a general way.\nAnother lighting source using an LED that emits blue light has been known from U.S. Pat. No. 6,791,259 B1. For producing white light, a phosphorous that emits yellow light, such as YAG:Ce or ZnS:Mn, is used, distributed in an epoxy resin or silicon layer or in a glass layer. In addition, a layer with radiation-diffusing particles is provided, which contains particles of BaTiO3, Al2O3, SiO2, CaCO3 or BaSO4. Those radiation-diffusing particles are again embedded in a carrier layer consisting of epoxy resin, silicon or glass.\nSuch a light-conversion material is connected with certain disadvantages due to the luminescent materials used being embedded in anorganic carrier materials. The granulates used give rise to scattering losses. And non-homogeneous distribution of the granulate on the solid-state emitter may lead to variable color perception as a function of angle. In addition, epoxy resins are instable over time in many respects, especially with respect to their optical and mechanical properties. And as a rule, their thermal stability is also unsatisfactory for generating high brightness. To the extent the respective diffusion materials and luminescence materials are embedded in glass carrier layers, the before-mentioned document does not provide any information how these are to be produced. Moreover, the production and application of a plurality of superimposed layers is complex and expensive.\nEP 1 642 869 A1 discloses a glass ceramic which preferably is used for down-conversion of excitation radiation in the blue and UV regions of the spectrum. The glass ceramic comprises the following components (on an oxide basis): 5-50% by weight of SiO2, 5-50% by weight of Al2O3, 10-80% by weight of Y2O3, 0-20% by weight of B2O3, 0.1-30% by weight of rare earths, preferably 15-35% by weight of SiO2, 15-35% by weight of Al2O3, 25-60% by weight of Y2O3, 1-15% by weight of B2O3 and 1-15% by weight of rare earths. The glass ceramic contains crystalline phases in which rare-earth ions are taken up at least in part. Crystalline phases containing yttrium ions as a component are replaced by rear-earth ions in this case at least in part. The phases in question may include, for example, Y3Al5O12 (YAG), Y2SiO5, Y2Si2O7, SrAl2O4, BaMgAl10O17, Sr2P2O7, Sr4Al14O25 and YbO3, that serve as host phase for taking up rare-earth ions at least in part.\nThe respective glasses are produced by a technological melting process and may then be ceramized. Ceramization is effected by initial tempering at a nucleation temperature of between 850° C. and 900° C., for a period of several hours, and then ceramizing at a temperature of between 1050 to 1150° C. for the time of one hour. The crystal phases identified in this case were Y2Si2O7, Y2SiO5, YbO3.\nThe conversion efficiency of such glass ceramics is, however, not yet sufficient for many applications because the glass ceramic contains a number of non-convertible crystal phases such as Y2Si2O7.\nJP(A) H04-119941 further discloses a glass ceramic that contains 5-50% by weight of SiO2, 5-70% by weight of Al2O3, 10-70% by weight of Y2O3 and 0.1 to 30% by weight of a nucleation agent such as MgO, TiO2, ZrO2 or La2O3. For the production process, starting materials (oxides) are mixed with organic solvents and binders, and are heated to then form a shaped glass by solid-state reaction. The glass so produced is then subjected first to a nucleation process, by tempering at temperatures of between 950° C. and 1010° C., and then to ceramization at a temperature of approximately 1100° C.\nThe production process is relatively complex. And in addition, the glass ceramic has a conversion efficiency not sufficient for all applications."}
{"text": "Super continuum (SC) generation is a nonlinear phenomenon characterised by dramatic spectral broadening of intense light passing through a nonlinear material. SC generation occurs in various media and may find use in numerous applications such as spectroscopy or ultra-short-pulse generation. Spectral slicing of a generated SC is an elegant way of replacing multiple light sources having separate wavelengths.\nMany prior art SC light sources are based on micro structured fibres (MF) which posses unique optical properties with a large degree of freedom of design thus allowing for optimization of the fibre for SC generation. However, SC generation in standard fibre has also been demonstrated. Most prior art disclosures of SC sources utilise femto-second (fs) pulses (10−15 s) to generate the SC. The physical mechanism responsible for the SC generation is believed to be due to multiple nonlinear processes. It has also been shown that it is possible to create SC by use of pico- and nanosecond pulses, and the mechanism responsible for these SCs is often attributed to a combination of four wave mixing and stimulated Raman scattering. In either case of pump duration and/or fibre type, the size of the guiding core has been relatively small (<6 μm diameter) to ensure a high optical density in the core so that non-linear effects are maximized. This limits the optical power with which the fibre may be pumped either due to the damage threshold of the end facet of the fibre or the damage threshold of the fibre core. The optical power may also be limited by the availability or cost of a high power laser having an output with a sufficient beam quality to allow a substantial amount of its output to enter the fibre."}
{"text": "1. Field of the Invention\nThe invention relates to a furniture assembly, more particularly to a furniture assembly, such as a chair, that is easy to assemble and that provides good support.\n2. Description of the Related Art\nGenerally, a conventional inflatable mattress or sofa bed relies only on air introduced into an inner portion thereof to maintain its shape and to support the weight of a user. When the air introduced into the conventional inflatable mattress or sofa bed is insufficient, the conventional inflatable mattress or sofa bed is uncomfortable and unsafe to use."}
{"text": "1. Field\nThe subject matter disclosed herein relates to an indoor likelihood heatmap and more specifically, but by way of example only, to generating or using an indoor likelihood heatmap in conjunction with movement within an indoor area.\n2. Information\nPaper maps have been used by people for hundreds, if not thousands of years, to aid navigation in unfamiliar or foreign territories. Electronic maps began to be available during the twentieth century. With the advent of the Internet, people could electronically access maps of many places from all over the globe. Web-based mapping services could also provide directions from one point to another point. These directions from web-based mapping services were originally relatively static. With the invention of satellite-positioning system (SPS) technology and ever-smaller electronic devices, however, so-called turn-by-turn directions could be provided dynamically as travelers journeyed toward their destination.\nElectronic maps, web-based mapping services, and turn-by-turn directions focus on providing navigational aids in certain situations and in particular environments. Unfortunately, there are other situations or different environments for which they are not intended or have not been designed. Consequently, there remain a number of situations, environments, etc. in which navigational or other location-based services may be improved."}
{"text": "Communication systems that support wireless and/or wired communications between wireless and/or wire-line communication devices are being rapidly proliferated. Such communication systems range from national and/or international cellular telephone systems, to the Internet, to point-to-point in-home wireless networks. Each type of communication system is typically constructed to operate in accordance with one or more communication standards. For instance, wireless communication systems may operate in accordance with one or more standards including, but not limited to, IEEE 802.11, Bluetooth, advanced mobile phone services (AMPS), digital AMPS, global system for mobile communications (GSM), code division multiple access (CDMA), local multi-point distribution systems (LMDS), and multi-channel-multi-point distribution systems (MMDS), among any number of other standards or combination of standards.\nDepending on the type of wireless communication system, a wireless communication device, such as a cellular telephone, two-way radio, personal digital assistant (PDA), personal computer (PC), laptop computer, home entertainment equipment, as some examples, may communicate directly and/or indirectly with other wireless or wired communication devices. For direct communication (e.g., point-to-point communications), the participating wireless communication devices may tune their receivers and transmitters to the same channel or multiple channels (e.g., one or more of the plurality of radio frequency (RF) carriers of the wireless communication system) and communicate over that channel or channels. For indirect wireless communications, each wireless communication device may communicate directly with an associated base station (e.g., for cellular services) and/or an associated access point (e.g., for an in-home or in-building wireless networks) via an assigned channel, or channels. To complete a communication connection between the wireless communication devices, the associated base stations and/or associated access points may communicate with each other directly, via a system controller, via the public switch telephone network, via the internet, and/or via some other wide area network.\nFor wireless communication devices to participate in wireless communications, the devices may each include a built-in radio transceiver (e.g., receiver and transmitter) or may be coupled to an associated radio transceiver (e.g., a station for in-home and/or in-building wireless communication networks, RF modem, etc.). As is known, the receiver portion of such a transceiver may receive RF signals, demodulate an RF carrier frequency from the RF signals via one or more intermediate frequency stages to produce baseband signals, and demodulate a baseband signal in accordance with a particular wireless communication standard to recapture transmitted data from an intermediate frequency signal. The transmitter portion of such a transceiver may convert data (e.g., digital data) into RF signals by modulating the data in accordance with a particular wireless communication standard to produce baseband signals and mix the baseband signals with an RF carrier using one or more intermediate frequency stages to produce RF signals. In other embodiments, transmitted and received signals may be directly converted between the RF domain and baseband.\nPhase locked loops (PLL) are often used in integrated wireless transceivers for RF frequency generation and RF signal modulation due to their high level of integration and ability to operate over a wide range of frequencies. As those working in this area are aware, two basic PLL topologies exist, which are generally referred to as “type 1” and “type 2,” respectively. A general structure is common to both types of PLLs, namely each topology includes a phase detector, a low-pass filter (LPF), a voltage controlled oscillator (VCO), and a feedback path. The phase detector serves as an “error amplifier” in the feedback loop. Based oh an “error” signal produced by the phase detector, the PLL may reduce any phase difference between the input reference signal and the feedback signal (e.g., “the error”). A PLL is considered “locked” if this phase difference is constant over time.\nIn “type 1” PLLs, the phase detector generates square voltage pulses whose duration is proportional to the phase error. These voltage pulses are filtered by the LPF to generate a smooth VCO control voltage whose amplitude is proportional to the phase error. The VCO responds to the change in control voltage by increasing or decreasing its oscillation frequency. The feedback action of the PLL then causes the VCO to lock to the desired operating frequency.\n“Type 2” PLLs, also known as “charge pump based PLLs,” include a phase and frequency detector, a charge pump, a loop filter, a voltage controlled oscillator (VCO), and a feedback path. The phase and frequency detector compares the phase and frequency of a reference signal with the phase and frequency of a feedback signal (e.g., the output oscillation produced by the VCO, which is fed back to the phase and frequency detector via the feedback loop). If the phase and/or frequency of the reference signal leads the phase and/or frequency of the feedback signal (which occurs when the output oscillation frequency is less than the desired oscillation frequency), the phase and frequency detector generates an up signal. In response to the up signal, the charge pump increases the positive current it outputs, which, when filtered by the loop filter, increases a control voltage applied to an input of the VCO. With an increase in the control voltage, the VCO increases its frequency of oscillation. If the phase and/or frequency of the reference signal lags the phase and/or frequency of the feedback signal (which occurs when the output oscillation frequency is greater than the desired oscillation frequency), the phase and frequency detector generates a down signal. In response to the down signal, the charge pump increases the negative current it outputs, which, when filtered by the loop filter, decreases the control voltage input to the VCO. With a decrease of the control voltage, the VCO decreases its frequency of oscillation.\nIn an “ideal” PLL (e.g., in a mathematical model of a PLL), the VCO operates linearly, which can be expressed as: Θout(t)=KVCO∫Vcntrl(t)dt, where KVCO is the gain (specified in MHz/volt, for example) of the VCO, Θout is the output oscillation of the VCO, and Vcntrl is the input control voltage of the VCO provided by the loop filter. Based on this mathematical model of PLL operation, when a PLL settles to a desired output frequency (i.e., the output signal frequency of the VCO), the input control voltage of the VCO assumes a constant value. When the PLL is used as a modulator, the control voltage may vary around this constant value in accordance with a desired modulation pattern, thus providing frequency modulation.\nIt follows from the above equation that the frequency of a modulated signal produced using a PLL operating in accordance with the mathematical model is proportional to the change in control voltage with a proportionality constant equal to KVCO. In practice, however, the gain of a VCO (i.e., KVCO) is not constant, but varies with integrated circuit fabrication process, PLL frequency band, control voltage amplitude, and temperature, among any number of other parameters. For instance, for a PLL that has a wide frequency range of operation, the gain of the VCO may vary from its nominal design value by as much as +/−50%. Since KVCO is a key parameter governing the dynamic behavior of the PLL, when used as a data signal modulator (e.g., in a translational loop), such large variations in VCO gain may produce prohibitively large modulation errors.\nOne possible solution to reduce the adverse affects of such variations in VCO gain is to design a baseband processor that is used in conjunction with the PLL to account for VCO gain variations. This approach, however, relies on an assumed VCO gain variation. Therefore, such an approach may not be particularly precise, making this solution of limited benefit."}
{"text": "Solid freeform fabrication (“SFF”) is the name given to a class of manufacturing methods which allow the fabrication of three-dimensional structures directly from computer-aided design (“CAD”) data. SFF processes are generally additive, in that material is selectively deposited to construct the product rather than removed from a block or billet. Most SFF processes are also layered, meaning that a geometrical description of the product to be produced is cut by a set of parallel surfaces (planar or curved), and the intersections of the product and each surface—referred to as slices or layers—are fabricated sequentially. Together, these two properties mean that SFF processes are subject to very different constraints than traditional material removal-based manufacturing. Nearly arbitrary product geometries are achievable, no tooling is required, mating parts and fully assembled mechanisms can be fabricated in a single step, and multiple materials can be combined, allowing functionally graded material properties. New features, parts, and even assembled components can be “grown” directly on already completed objects, suggesting the possibility of using SFF for the repair and physical adaptation of existing products. On the other hand, a deposition process must be developed and tuned for each material, geometry is limited by the ability of the deposited material to support itself and by the (often poor) resolution and accuracy of the process, and multiple material and process interactions must be understood.\nSFF has traditionally focused on printing passive mechanical parts or products in a single material, and the emphasis of research has been on developing new deposition processes (U.S. Pat. No. 5,121,329 to Crump; U.S. Pat. No. 5,134,569 to Masters; U.S. Pat. No. 5,204,055 to Sachs et. al.; and U.S. Pat. No. 5,126,529 to Weiss et. al.), on improving the quality, resolution, and surface finish of fabricated products, and on broadening the range of single materials which can be employed by a given SFF process, including biocompatible polymers and other biomaterials (Pfister et al., “Bio functional Rapid Prototyping for Tissue-engineering Applications: 3D Bioplotting Versus 3D Printing,” Journal of Polymer Science Part A: Polymer Chemistry 42:624-638 (2004); Landers et al., “Desktop Manufacturing of Complex Objects, Prototypes and Biomedical Scaffolds by Means of Computer-assisted Design Combined with Computer-guided 3D Plotting of Polymers and Reactive Oligomers,” Macromolecular Materials and Engineering 282:17-21 (2000)), and living cells (Roth et al., “Inkjet Printing for High-throughput Cell Patterning,” Biomaterials 25:3707-3715 (2004)). These improvements have allowed freeform fabrication to become a viable means of manufacturing finished functional parts, rather than only prototypes.\nMore recently, the greater utility of freeform fabricated products having multiple materials has been recognized, prompting reexamination and novel research into processes which can fabricate using multiple materials (U.S. Pat. No. 5,260,009 to Penn), and which can thereby produce complex articles with a variety of functionality, including functionally graded materials (Ouyang et al., “Rapid Prototyping and Characterization of a WC-(NiSiB Alloy) Ceramet/Tool Steel Functionally Graded Material (FGM) Synthesized by Laser Cladding,” Columbus, Ohio, USA: TMS—Miner. Metals & Mater. Soc. (2002); Smurov et al., “Laser-assisted Direct Manufacturing of Functionally Graded 3D Objects by Coaxial Powder Injection,” Proceedings of the SPIE—The International Society for Optical Engineering 5399:27 (2004)), electronics, MEMS (Fuller et al., “Ink-jet Printed Nanoparticle Microelectromechanical Systems,” Journal of Microelectromechanical Systems 11:54-60 (2002)), living tissue constructs (Mironov et al., “Organ Printing: Computer-aided Jet-based 3D Tissue Engineering,” Trends in Biotechnology 21:157-161 (2003)), and compositions of living and nonliving materials (Sun et al., “Multinozzle Biopolymer Deposition for Tissue Engineering Application,” 6th International Conference on Tissue Engineering, Orlando, Fla. (Oct. 10-13, 2003); International Patent Application No. PCT/US2004/015316 to Sun et al.; and U.S. Pat. No. 6,905,738 to Ringeisen et al.). All of these systems still depend upon a small fixed set of deposition process technologies, and are therefore limited to the materials which can be adapted to those processes, by the effects of those particular processes on the materials, and by the fabrication rates and resolutions of those processes. In particular, the system of U.S. Pat. No. 6,905,738 to Ringeisen et al., requires that for every material to be deposited, a two material system be developed comprising the material to be transferred, and a compatible matrix material which is vaporized by the laser in order to propel the transfer material to the substrate. In addition, this system has only demonstrated fabrication of thin films of materials—its ability to deposit many layers of materials is not well established. The system and method of Sun et al., “Multinozzle Biopolymer Deposition for Tissue Engineering Application,” 6th International Conference on Tissue Engineering, Orlando, Fla. (Oct. 10-13, 2003) and International Patent Application No. PCT/US2004/015316 to Sun et al., is limited to a fixed set of four deposition processes and requires that the alginate materials be deposited into a bath of liquid crosslinking agent—a limitation it shares with the work of Pfister et al., “Bio functional Rapid Prototyping for Tissue-engineering Applications: 3D Bioplotting Versus 3D Printing,” Journal of Polymer Science Part A: Polymer Chemistry 42:624-638 (2004) and Landers et al., “Desktop Manufacturing of Complex Objects, Prototypes and Biomedical Scaffolds by Means of Computer-assisted Design Combined with Computer-guided 3D Plotting of Polymers and Reactive Oligomers,” Macromolecular Materials and Engineering 282:17-21 (2000). In addition, none of these systems explicitly measures the properties of, and monitors and controls the conditions experienced by the fabrication materials, the fabrication substrate, and the article under construction before, during, and/or after fabrication as an intrinsic part of the fabrication process and manufacturing plan. The fabrication process is thus limited to the spatial control of material placement on relatively simple, passive substrates. Temporal control of the evolution of material properties is therefore not possible, and complex substrates whose state must be controlled and monitored continuously are not readily accommodated. Fabricating into or onto substrates, such as living organisms or devices which must remain in operation continuously, is problematic.\nA major challenge in orthopaedic tissue engineering is the generation of seeded implants with structures that mimic native tissue, both in terms of anatomic geometries and intratissue cell distributions. Previous studies have demonstrated that techniques such as injection molding (Chang et al., “Injection Molding of Chondrocyte/Alginate Constructs in the Shape of Facial Implants,” J. Biomed. Mat. Res. 55:503-511 (2001)) and casting (Hung et al., “Anatomically Shaped Osteochondral Constructs for Articular Cartilage Repair,” J. Biomech. 36:1853-1864 (2003)) of hydrogels can generate cartilage tissue in complex geometries. Other studies have investigated methods to reproduce regional variations in articular cartilage constructs by depositing multiple layers of chondrocytes (Klein et al., “Tissue Engineering of Stratified Articular Cartilage from Chondrocyte Subpopulations,” Osteoarthritis Cartilage 11:595-602 (2003)) or chondrocyte-seeded gels (Kim et al., “Experimental Model for Cartilage Tissue Engineering to Regenerate the Zonal Organization of Articular Cartilage,” Osteoarthritis Cartilage 11:653-664 (2003)). However, there remains no viable strategy for rapidly producing implants with correct anatomic geometries and cell distributions. Recently, advances in SFF techniques have enabled the deposition of multilayered structures composed of multiple chemically active materials (Malone et al., “Freeform Fabrication of 3D Zinc-Air Batteries and Functional Electro-Mechanical Assemblies,” Rapid Prototyping Journal 10:58-69 (2004)). Applicants believe that this technology has the potential to be adapted to the fabrication of living tissue under conditions that preserve cell viability.\nThe present invention is directed at overcoming disadvantages of prior art approaches and satisfying the need to establish a robust and reliable SFF system and method."}
{"text": "1. Field of the Invention\nThe invention relates to a circuit arrangement for the transmission of telecommunication signals between a switching center of a telecommunication network and at least one subscriber.\n2. Description of the Prior Art\nIn telecommunications networks, subscribers are connected to a switching center by means of a connecting line. At the subscriber's end of the connecting line, there is connected a device in which, in addition to other passive and active internal circuit elements, a direct-current source and a transformer with a core, an input winding, and an output winding are arranged. The input winding is connected with the connecting line via the internal circuit elements of the device. The output winding is connected to the subscriber by means of another connecting line. The output winding of the transformer is subdivided into two partial windings of identical size, between which there are connected a first capacitor and a switch connected in series with the windings. The switch is to be actuated by a component for detecting the loop current. The transformer is equipped with at least one additional winding, whose number of turns is equal to the total of the number of turns of the two partial windings. In the additional winding, the current flows in a direction opposite to that in which it flows in the two partial windings. The direct-current source is connected with one of the partial windings at the connection point of the first capacitor, with the interposition of the additional winding, as shown in German OS 37 15 552.\nA circuit arrangement of this type is used, for example, in the so-called \"PCM 2\" process, in which two subscribers in adjacent locations of a telecommunication network are jointly connected to a switching center by means of a connecting line. The letters \"PCM\" stand for the known pulse code modulation process, by means of which telecommunication signals are transmitted as digital signals. The additional explanations apply to this process, representing all possible other cases of application of the invention.\nIn the PCM 2 process, analog signals are digitalized in the switching center and fed into the connecting line in coded form. At the end of the connecting line there is installed a device in which the incoming digital signals are decoded and transformed back into analog signals. This, of course, applies also to the opposite direction. The two subscribers are connected to the device by means of connecting lines. For each subscriber, the device contains a transformer whose input winding is connected to the connecting line with the interposition of an analog-digital converter and other circuit elements. Each particular subscriber is connected to the output winding of the transformer by means of a connecting line. Each subscriber is assigned his own direct-current source for supplying his telephone, which is connected at the output side of the transformer. Particularly in the cases where the subscribers are connected to the device by means of long connecting lines, interferences can be induced by external electric and magnetic fields, which can lead to cross-talk in other channels and to humming sounds.\nTo prevent these interferences, in the known circuit arrangement according to the previously mentioned German OS 37 15 552, the output winding of the transformer is subdivided into two partial windings of equal size, between which the first capacitor and the switch are connected in series. By dividing the output winding into two partial windings, the whole arrangement is made symmetrical in such a way that identical conditions always result for the forward and return lines of the connecting line with respect to a reference potential--for example, the ground. Interferences induced by electrical or magnetic fields can therefore no longer have an effect, regardless of the length of the connecting lines. The first capacitor makes certain that the direct current supplied flows only in the desired direction.\nThe brief interruptions in the feed current occurring during a dialing process cannot be rapidly determined by the component detecting the loop current, because current initially continues to flow through the first capacitor. It may therefore happen that the switch, depending on the capacity of the first capacitor, is opened with such a delay that a distortion of the dialing characters takes place. This can lead to wrong connections."}
{"text": "A ventriculoperitoneal (VP) shunt is used in Hydrocephalus cases to drain excess cerebral spinal fluid (CSF) from the ventricles to reduce the pressure in the ventricular cavity. Some VP shunts may be equipped with an adjustable or programmable shunt valve configured to open and drain excess CSF in an instance in which a predetermined pressure set point is reached within the ventricular cavity.\nIn some cases, the VP shunt may become blocked or clogged preventing pressure relief in the ventricular cavity; however, the symptoms of a blocked VP shunt and the symptoms of nasal congestion, typical of a common cold, may be substantially similar. A ventricular pressure measurement is used to differentiate between nasal congestion and blockage of the VP shunt. In some cases, invasive procedures such as drilling a hole in the skull are used to take a ventricular pressure measurement. In other instances, computerized axial tomography (CT) scans or X-rays may be taken over a period of several days and changes in the CT or X-ray images may be indicative of a ventricular pressure change. In still another instances, an ultrasound of the retinal vessel may be used for indirect measurement of ventricular pressure. As such, there is no non-invasive method to quickly and directly measure ventricular pressure."}
{"text": "1. Technical Field\nThe present disclosure relates to orthopaedic prostheses and, specifically, to articular tibial components in a knee prosthesis.\n2. Description of the Related Art\nOrthopaedic prostheses are commonly utilized to repair and/or replace damaged bone and tissue in the human body. For a damaged knee, a knee prosthesis may be implanted using a tibial baseplate, a tibial bearing component, and a distal femoral component. The tibial baseplate is affixed to a proximal end of the patient's tibia, which is typically resected to accept the baseplate. The femoral component is implanted on a distal end of the patient's femur, which is also typically resected to accept the femoral component. The tibial bearing component is placed between the tibial baseplate and femoral component, and may be fixed upon or slidably coupled to the tibial baseplate.\nThe tibial bearing component, which may also be referred to as a tibial insert or meniscal component, provides an articular surface which interacts with the adjacent femur or femoral component during extension and flexion of the knee. The features and geometry of the articular surface influences the articular characteristics of the knee, such as by defining maximum knee flexion, internal/external rotation, femoral rollback, and behavior of the knee prosthesis in hyperextension, for example. Accordingly, substantial design efforts have previously focused on providing knee prosthesis components which preserve flexion range and promote a desired kinematic motion profile for the widest possible range of prospective knee replacement patients."}
{"text": "1. Technical Field\nThe present disclosure relates to positioning assemblies, and particularly, to a positioning assembly having a fastener.\n2. Description of the Related Art\nElectronic devices such as flat panel TVs often defines a mounting hole in a rear face to enable wall mounting. When the electronic device is to be mounted, a threaded member is partially received in a threaded hole in the wall and also partially received in the mounting hole in the rear face. However, the threaded member is prone to detachment since the allowable length thereof is limited.\nTherefore, there is room for improvement within the art."}
{"text": "When surfing online with a browser operating on a computer or smart phone, users often navigate through web pages and wish to share with others some pages they like. For example, when a user is browsing a web page, she/he may copy the URL in the browser address bar, paste the copied URL link to an instant messaging window (e.g., MSN or QQ) to share the page with a friend chatting in the messaging window in real time. After receiving the URL, the friend may click on the link pasted in the window, and browse the page.\nThis example of web page sharing requires switching between the browser and the instant messaging tool, which means the user must have installed the instant messaging tool in addition to the browser."}
{"text": "Exemplary embodiments of the present invention relate generally to games and more particularly to card games in which players attempt to discard all of the cards in their hand to win the game. Various embodiments also relate generally to a device that can be used for playing a card game and methods of playing games using the device. The play of various card games is facilitated by one or more cards and the game playing device.\nExamples of cards, card games, and devices for card games are disclosed in U.S. Pat. No. 5,114,153 entitled “Mechanical Card Dispenser and Method of Playing a Card Game,” issued May 19, 1992, and in U.S. Patent Application Publication No. 2007/0108703A1, entitled “Card and Selection Device Games,” published May 17, 2007, the disclosure of each of which is incorporated by reference herein in its entirety.\nConventional games exist in which players attempt to be the first one to dispose of all of their cards, or try not to be the last player still holding one or more cards when the other players do not have any cards. An example of such a game is the UNO® card game in which each player strives to be the first one to dispose of all of the cards from their hand. The playing cards from the UNO® card game are owned by Mattel, Inc. The UNO® card game involves a discard pile of cards, a draw pile of cards and cards that are dealt to each player. In the UNO® card game, a player not being able to match the upwardly facing card in a discard pile by number or by color, and not having a wild card, must draw cards from a draw pile until the player is able to make the required match or play a wild card."}
{"text": "This application claims the benefit of French application no. 00/10088 filed Jul. 31, 2000.\nThe present invention relates to a method of securing access to at least one application supported by a user card having a microprocessor. The user card co-operates with one of the communications terminals in a communications system. The invention also relates to a communications terminal implementing the method.\nA non-exclusive field of application of the invention is that of mobile terminals for radio communication operating in a cellular radio communication system. The invention applies particularly but not exclusively to a system implementing the global system for mobile communication (GSM).\nThe terminals of a radio communication system can co-operate with at least one smart card. Each smart card contains at least one internal application. Each internal application is suitable for co-operating with an application external to the smart card. For example, the external application can be a bank application co-operating with a user card of the payment card type supporting the corresponding bank application.\nIn addition, when an external application seeks to access the corresponding internal application supported by the smart card, said external application starts the connection protocol. It exchanges so-called applicative protocol data unit (APDU) data via the application programming interface (API) of the terminal. In the APDU, the external application specifies its class or type (bank, etc.). This class is given in one of the bytes of the APDU which is referred to as the xe2x80x9cCLAxe2x80x9d byte.\nThe concept of having a plurality of applications united on a common card is highly advantageous for the subscriber. This enables the subscriber to use the terminal to perform numerous operations such as paying for an order likewise made from the terminal, and to do so simply and solely by means of the terminal.\nThus, various different external applications can have access simultaneously to corresponding internal applications on a single smart card. However, two identical external applications must not have access simultaneously to the corresponding internal application on the smart card.\nIt must be possible to restrict access in this way in order to mitigate the risks of fraudulent use of smart card applications. During a transaction between an external bank application and the bank application on a subscriber\"\"s smart card, it is necessary to avoid any transaction with a pirate application of the bank type requesting access to the subscriber\"\"s smart card.\nFurthermore, within a smart card containing a plurality of different internal applications, it is difficult to guarantee that there will be no leakage between the various internal applications. Thus, consider the case of a first external application X already in communication with the corresponding internal application X situated on the smart card. If a pirate external application Xxe2x80x2 seeks to communicate with the internal application X, then the application Xxe2x80x2 might be able to connect with an internal application Y situated in the smart card close to the internal application X. Thereafter, the application Xxe2x80x2 might be able to take advantage of any possible leakage between the internal applications in order to reach the internal application X.\nIn the state of the art, there exists an on/off model. That model operates, for example, on pieces of equipment such as Internet communication terminals that give access to the Internet. With this type of terminal, applications can be downloaded from the outside.\nIf the connection protocol used by the external application is known by the terminal, then the external application is said to be an xe2x80x9cauthorizedxe2x80x9d application and it can access the corresponding internal application of the smart card. However, if another authorized external application of the same class as the first requests access to the same internal application of the smart card, the second application can also gain access to the same internal application of the smart card, thereby establishing an interface between the two external applications.\nThe prior solution is therefore not sufficient for mitigating the above-mentioned risk.\nThe invention solves this problem by enabling an external application to reserve access to the corresponding internal application of the user card to itself alone.\nMore precisely, the invention provides a method of securing access of an external application to at least one user card having a microprocessor.\nSaid user card is capable of containing a plurality of internal applications, with said external application corresponding to one of the internal applications.\nSaid user card co-operates with a terminal.\nSaid external application uses a protocol for connection to said user card in which each type of external application is associated with a predetermined identify parameter. The method comprises the following steps:\nwhen a first external application is connected with the corresponding internal application on the user card, the identity parameter of the application is determined;\nif a second external application requests connection to said user card, the identity parameter of said second application is analyzed; and\nif said identity parameter is identical to that of the first-connected application, then access to said corresponding internal application on the user card is denied to said second external application.\nIn particular, the method comprises a connection protocol between the external application and the corresponding internal application on the user card. This protocol results in an APDU exchange, containing a plurality of bytes, between the external application and the internal application situated on the user card via the software interface of the terminal, said identity parameter being represented by at least one of the bytes of the APDU.\nThe method comprises a connection protocol during which an APDU exchange takes place such that at least one of the bytes of the APDU is the CLA byte defining the class of the external application, each identical external application type that might be capable of exchanging APDUs having an identical CLA byte, said identity parameter being the CLA byte.\nThe method is such that once the external application of a defined class has fully reserved access to at least one user card, said identity parameter is analyzed by the software interface for all possible classes which define an external application.\nOther characteristics and advantages of the invention appear on reading the following description relating to an embodiment given by way of non-limiting example. The description refers to the accompanying drawings."}
{"text": "1. Technical Field\nThe present disclosure relates to acetabular implants for hip replacement surgery. In particular, the present disclosure relates to acetabular implants including a shell component and alternative bearing components interchangeably engageable with the shell component to provide a choice of bearing materials.\n2. Description of the Related Art\nTotal hip replacement surgery is commonly performed to alleviate pain and loss of function in injured or diseased hip joints. During this surgery, the articulating surfaces of the hip joint are replaced with prosthetic bearing components. The replacement components generally include a femoral component having a convex bearing surface and an acetabular cup component having a mating concave bearing surface.\nModular prosthetic components have become popular because they allow the surgeon to assemble components in a variety of configurations at the time of surgery to meet specific patient needs and surgeon preferences. For example, modular acetabular components generally include separate shell and liner components that can be assembled in a variety of configurations of shell surface finish, shell outer diameter, liner inner diameter, and liner bearing material. With a modular acetabular component, it is desirable to lock the shell and liner together to prevent expulsion of the liner and to minimize debris producing wear between them. Typically, the engagement mechanism is formed adjacent the equator of a hemispherical shell and liner to maximize the engagement area and the resulting holding power of the engagement mechanism.\nVarious liner bearing materials are in use. The liners vary in hardness, friction coefficient with different paired ball heads, weight, and wear resistance. Polymers, including ultrahigh molecular weight polyethylene (UHMWPE), are commonly used as bearing materials paired with an opposing metal, ceramic, or other composition ball head. The wear resistance of UHMWPE has been improved by irradiating it to cause changes in its chemical and mechanical properties. As the wear properties are improved, the bulk physical properties also change. Other materials, including metals and ceramics, have also been used for acetabular bearings. These materials vary from one another in terms of their hardness, resilience, brittleness, and other physical properties. Because of this variation, various mechanisms have been developed for engaging acetabular liners with their mating shells. Different engagement mechanisms are suitable for different liner and shell material combinations. These engagement mechanisms include snap-fit, cylindrical press-fit, taper-fit, threaded engagement, and other suitable locking mechanisms. It is desirable to be able to alternately fit different liners into a common shell to reduce inventory while allowing surgeon choice in liner selection. It is also desirable to allow intraoperatively changing from one liner to another without having to remove a shell that has already been placed in the surgical site during a primary surgery, or one that has become well fixed and only needs liner replacement in a revision surgery.\nU.S. Pat. No. 6,475,243 issued to Sheldon et al. on Nov. 5, 2002 (the “'243 patent”). The '243 patent discloses an acetabular cup assembly that allows pre-operative or inter-operative selection and securement of a bearing member within a shell member. The shell is metallic, while the bearing insert is in the form of a plastic bearing member. The assembly of the '243 patent includes a securing mechanism including an annular recess formed in the shell and a complementary annular rib that seats within the annular recess. In this assembly, the preferred material for the bearing member is polyethylene.\nAlternatively, the '243 patent discloses a sleeve that may be secured within the shell by locking of tapered seating surfaces. A bearing member is secured to the sleeve with a tapered securement surface. The preferred material for the sleeve is commercially pure titanium. A recess may be formed in the tapered seating surface so that engagement of tapered seating surfaces occurs along two segments having generally the same axial length.\nU.S. Pat. No. 6,610,097 issued to Serbousek et al. on Aug. 26, 2003 (the “'097 patent”). The '097 patent discloses a shell, a liner and a bearing. The bearing is coupled to the liner to form a subassembly by cooling the bearing, such as with liquid nitrogen, to shrink the size of the bearing, then press-fitting the bearing into the liner and allowing the bearing to warm and return to a larger size. The bearing and liner are fastened together in a fixed and locked position to form the bearing/liner subassembly. The liner provides a metal taper surface that forms a metal-to-metal locking connection between the subassembly and the shell."}
{"text": "1. Field of the Invention\nThe invention relates to optical scanners for bar code reading, and in particular to the decoding algorithm used in a laser scanning bar code reader for processing blur or reducing the noise in reflected light signals.\n2. Description of the Related Art\nBar code symbols are formed from bars or elements typically rectangular in shape with a variety of possible widths. The specific arrangement of elements defines the character represented according to a set of rules and definitions specified by the code or “symbology” used. The relative size of the bars and spaces is determined by the type of coding used as is the actual size of the bars and spaces. The number of characters (represented by the bar code symbol) is referred to as the density of the symbol. To encode the desired sequence of the characters, a collection of element arrangements are concatenated together to form the complete bar code symbol, with each character of the message being represented by its own corresponding group of elements. In some symbologies, a unique “start” and “stop” character is used to indicate when the bar code begins and ends. A number of different bar code symbologies are in widespread use including UPC/EAN, Code 39, Code 128, Codeabar, and Interleaved 2 of 5.\nIn order to increase the amount of data that can be represented or stored on a given amount of surface area, several more compact bar code symbologies have been developed. One of these code standards, Code 49, exemplifies a “two dimensional” symbol by reducing the vertical height of a one-dimensional symbol, and then stacking distinct rows of such one dimensional symbols, so that information is encoded both vertically as well as horizontally. That is, in Code 49, there are several rows of bar and space patterns, instead of only one row as in a “one dimensional” symbol. The structure of Code 49 is described in U.S. Pat. No. 4,794,239. Another two-dimensional symbology, known as “PDF417”, is described in U.S. Pat. No. 5,304,786.\nStill other symbologies have been developed in which the symbol is comprised not of stacked rows, but a matrix array made up of hexagonal, square, polygonal and/or other geometric shapes, lines, or dots. Such symbols are described in, for example, U.S. Pat. Nos. 5,2276,315 and 4,794,239. Such matrix code symbologies may include Vericode, Datacode, and MAXICODE.\nVarious optical scanning systems and readers have been developed heretofore for reading indicia such as bar code symbols appearing on a label or on the surface of an article. The bar code symbol itself is a coded pattern of graphic indicia comprised of a series of bars of various widths spaced apart from one another to bound spaces of various widths, the bars and spaces having different light reflecting characteristics. The readers function by electro-optically transforming the spatial pattern represented by the graphic indicia into a time-varying electrical signal, which is in turn decoded into data which represent the information or characters encoded in the indicia that are intended to be descriptive of the article or some characteristic thereof. Such data is typically represented in digital form and utilized as an input to a data processing system for applications in point-of-sale processing, inventory control distribution, transportation and logistics, and the like. Scanning systems and readers of this general type have been disclosed for example, in U.S. Pat. Nos. 4,251,798; 4,369,361; 4,387,297; 4,409,470; 4,760,248; 4,896,026; 5,015,833; 5,262,627; 5,504,316; and 5,625,483, all of which have been assigned to the same assignee as the instant application and each of which are hereby incorporated by reference. As disclosed in some of the above patents, one embodiment of such a scanning system resides, inter alia, in a hand held, portable laser scanning device supported by a user, which is configured to allow the user to aim the scanning head of the device, and more particularly, a light beam, at a targeted symbol to be read.\nThe light source in a laser scanner bar code reader is typically a semiconductor laser. The use of semiconductor devices as the light source is especially desirable because of their small size, low cost and low voltage requirements. The laser beam is optically modified, typically by an optical assembly, to form a beam spot of a certain size at the target distance. It is often preferred that the cross section of the beam spot measured in the scanning direction at the target distance be approximately the same as the minimum width in the scanning direction between regions of different light reflectivity, i.e., the bars and spaces of the symbol. Although typical readers utilize a single laser source, other bar code readers have been proposed with two or more light sources of different characteristics, e.g. different frequencies.\nTypically, the semiconductor lasers used in such bar code scanner is an edge-emitting injection laser in which the laser beam in emitted from the p-n junction region on a polished end face of the device. By their physical nature, these known edge-emitting injection lasers emit a beam from a thin region at the p-n junction. A laser beam emanating from a thin source has a large beam divergence which makes focusing difficult and results in a wide rang of variability in performance from laser to laser.\nA more recently developed from of semiconductor laser is the vertical-cavity surface-emitting laser diode (VCSEL), such as described in “Efficient Room-Temperature Continuous-Wave AlGaInP/Al GaAs Visible (670 nm) Vertical-Cavity Surface Emitting Laser Diodes” by R P Schneider et al. published IEEE Photonics Technology Letters, Vol. 6, No. 3, Mar. 1994. Reference is also made to U.S. Pat. Nos. 5,283,447; 5,285,455; 5,266,794; 5,319,496; and 5,326,386, which are hereby incorporated by reference, for background information.\nThe VCSEL has substantial surface area from which the laser beam is emitted; this are may be patterned. Thus, the beam produced is less divergent in one dimension than with know edge-emitting type semiconductor laser diodes. The output beam is round, and is virtually not astigmatic. Furthermore, VCSELs typically operate at significantly lower current than edge-emitting laser diodes. Therefore, it also generates less heat.\nIn the laser beam scanning systems known in the art, a single laser light beam is directed by a lens or other optical components along the light path toward a target that includes a bar code symbol on the surface. The moving-beam scanner operates by repetitively scanning the light beam in a line or series of lines across the symbol by means of motion or a scanning component, such as the light source itself or a mirror disposed in the path of the light beam. The scanning component may either sweep the beam spot across the symbol and trace a scan line across the pattern of the symbol, or scan the field of view of the scanner, or do both. The laser beam may be moved by optical or opto-mechanical means to produce a scanning light beam. Such action may be performed be either deflecting the beam (such as by a moving optical element, such as a mirror) or moving the light source itself. U.S. Pat. No. 5,486,944 describes a scanning module in which a mirror is mounted on a flex element for reciprocal oscillation by electromagnetic actuation. U.S. Pat. No., 5,144,120 to Krichever et al., describes laser, optical and sensor components mounted on a drive for repetitive reciprocating motion either about an axis or in plane to effect scanning of the laser beam.\nAnother type of bar code scanner employs electronic means for causing the light beam to be deflected and thereby scan a bar code symbol, rather than using a mechanical motion to move or deflect the beam. For example, a linear array of closely spaced light sources activated one at a time in a regular sequence may be transmitted to the bar code symbol to simulate a scanned beam for a single source. Instead of a single linear array of light sources, a multiple-line array may also be employed, producing multiple scan lines. Such type of bar code reader is disclosed in U.S. Pat. No. 5,258,605 to Metlitsky et al.\nBar code reading systems also include a sensor or photodetector which detects light reflected or scattered from the symbol. The photodetector or sensor is positioned in the scanner in an optical path so that it has a field of view which ensures the capture of a portion of the light which is reflected or scattered off the symbol, detected, and converted into an electrical signal.\nIn retroreflective light collection, a single optical component e.g., a reciprocally oscillatory mirror, such as described in Krichever et al. U.S. Pat. No. 4,816,661 or Shepard et al. U.S. Pat. No. 4,409,470, both herein incorporated by reference, and Ser. No. 08/727,944, filed Oct. 9, 1996, scans the beam across a target surface and directs the collected light to a detector. The mirror surface usually is relatively large to receive as much incoming light as is possible, only a small detector is required since the mirror can focus the light onto a small detector surface, which increases signal-to-noise ratio.\nOf course, small scan elements are preferable because of the reduced energy consumption and increased frequency response. When the scan element becomes sufficiently small, however, the area of the scanning mirror can no longer be used as the aperture for the received light. One solution is to use a staring detection system (a non-retroreflective system) which receives a light signal from the entire field which the scanned laser spot covers.\nIn non-retroreflective light collection, the reflected laser light is not collected by the same optical component used for scanning. Instead, the detector is independent of the scanning beam, and is typically constructed to have a large field of view so that the reflected laser light traces across the surface of the detector. Because the scanning optical component, such as a rotating mirror, need only handle the outgoing light beam, it can be made much smaller. On the other hand, the detector must be relatively large in order to receive the incoming light beam from all locations in the scanned field.\nElectronic circuitry and software decode the electrical signal into a digital representation of the data represented by the symbol that has been scanned. For example, the analog electrical signal generated by the photodetector may be converted by a digitizer into a pulse width modulated digitized signal, with the widths corresponding to the physical widths of the bars and spaces. Alternatively, the analog electrical signal may be processed directly by a software decoder. See, for example, U.S. Pat. No. 5,504,318.\nThe decoding process of bar code reading systems usually works in the following way. The analog signal from the sensor or photodetector may initially be filtered and processed by circuitry and/or software to remove noise, adjust the dynamic range, or compensate for signal non-uniformities. Samples may then be taken of the analog signal, and applied to an analog-to-digital converter to convert the samples to digital data. See, for example, U.S. Pat. No. 6,170,749, which is herby incorporated by reference. Alternatively, analog circuitry may be used to digitize the shape of the signal. The digital samples or the pulse width modulated digitized signal is applied to a software algorithm, which attempts to decode the data in the signal to characters in the symbology. If the characters of the bar code symbol are all decoded successfully, the decoding process terminates and an indicator of a successful read (such as a green light and/or audible beep) is provided to the user. Otherwise, after a period of time, the decoder aborts operation of the decoding algorithm on the first scan data, receives the next scan, and performs another decode operation, and so on, until the scan data is decoded according to symbology specification into a binary representation of the data originally encoded in the symbol, and to the alphanumeric characters so represented.\nThe binary data representing the decoded characters is communicated to a host computer by an interface cable or wireless communication link. The interface cable may be a “smart cable” such as that described in U.S. Pat. Nos. 5,664,229 and 5,675,139, the contents of which are hereby incorporated by reference."}
{"text": "This disclosure relates to optical alignment systems and more particularly, to a system and method for improving an alignment routine for lithography or pattern recognition.\nIntegrated circuit chips are fabricated one level at a time. The levels include diffusions, gates, metal lines, insulation, isolation, and contacts. The structures on these levels must be precisely positioned so that the finished chip has structures properly positioned. The step of positioning a level with respect to a previously formed level is called alignment.\nCurrent industry methodologies require alignments of some feature to allow for orientation of a part. More specifically, alignment of patterned materials involves manual selection of alignment targets. The alignment targets typically include two or three alignment marks which a recognition system uses to learn a position of the part or material subject to examination. The user then provides a region of interest (ROI) by manually moving to this location while the software records this coordinate. Alternatively, a coordinate relative to some known reference point (e.g., center of the part) is indicated by the software.\nFurthermore, if multiple patterns exist on the part as found in semiconductor wafers, then a step periodicity is supplied to find the next ROI. Once the setup is completed (i.e., after the alignment marks are recorded) the optical system aligns and moves to a ROI on the part for measurement or inspection.\nOne drawback to the above approach is that time is needed to teach the alignment mark locations relative to a ROI and is wasted time. Furthermore, pattern recognition systems which fail to align are unable to re-teach themselves since no point of origin has been established."}
{"text": "This invention relates to a system and method for completing at least two calls from a premises that is connected by a single point-of-attachment having a single phone number to a telephone network and more particularly a public switched telephone network using an analog interface.\nMultiple calls often are established over digital networks from terminals connected by a point-of-attachment identified by a single terminal number. For example, networks using packet-switched data networks, described in ITU-T recommendations X.25, use digital terminations and allow multiple calls over the network. Other examples include frame-relay networks, described in ITU-T recommendation I.233; narrow-band ISDN, described in ITU-T recommendation 1.210; and broad-band ISDN networks, described in ITU-T recommendation I.211. These services use digital terminations that can be expensive to implement and maintain.\nIt would be advantageous if a premises connected by a point-of-attachment to the public switched Telephone Network (PSTN) having a single phone number for handling voice calls with the standard, analog xe2x80x9ctip-and-ringxe2x80x9d interface could carry a modulated signal between the premises and the public switched telephone network that includes multi-call capability from terminal equipment located on the premises, such as a telephone, video, fax machine, or computer. This would be advantageous, for example, if someone at a premises is using a computer while at the same time a telephone call needs to be made to a different location. Other examples could include someone making two telephone calls to the network, or using a computer and fax machine or a telephone and fax machine simultaneously.\nIt is therefore desirable if a method and system could complete at least two calls to different locations originating from terminal equipment located at a premises, which is connected by a point-of-attachment having a single phone number to a public switched telephone network without digital termination such as ISDN systems.\nThe present invention now allows a method and system for completing at least two calls originating from terminal equipment located at a premises and connected by a point-of-attachment having a single phone number to a telephone network where one dial tone is provided on an analog interface to the premises. The method and system of this invention permits multiple calls, such as an analog phone, video, data, fax and comparable calls over a single analog interface. It can be extended to allow an increased number of calls by utilizing additional analog interfaces. Such an extension, sometimes known as a multi-link arrangement, provides for any allocation of the multiple calls across the plurality of analog interfaces.\nIn accordance with one aspect of the present invention, a control channel is established between the telephone network, such as a public switched telephone network, and the terminal equipment located at the premises. At least two calls originating from premises equipment and the control channel are combined at the premises into a modulated analog signal having a separate channel for each call. Likewise, any incoming calls are combined with the above calls and the control channel for transmission between the customer premises and the public switched telephone network. Identity information is transmitted from the originating terminal equipment through the control channel to aid in determining a remote device with which a call is desired. This includes auxiliary information for progressing the call. The signal is transmitted from a modulation/demodulation device as a modulated analog signal to the public switched telephone network. The modulated analog signal is split into respective calls and the control channel, and the calls are then forwarded to a respective remote device based on identity information and auxiliary information received through the control channel.\nIn one aspect of the present invention, the calls can be split at the public switched telephone network. The public switched telephone network can also include a cellular network. Identifiers are associated with each channel for distinguishing each channel within the public switched telephone network. Each channel can also comprise a plurality of information blocks, and identifiers can be associated between the public switched telephone network and the terminal equipment for each channel to distinguish the information blocks belonging to each call or control channel."}
{"text": "1. Technical Field\nThe present invention pertains to the field of accessories for automotive vehicles. More particularly, the present invention pertains to a front bumper for an automotive vehicle.\n2. Discussion of Related Art\nAn ordinary front bumper is a shield against low-speed collisions, typically made of steel, aluminum, rubber, or plastic, and mounted on the front of a vehicle. When a low-speed collision occurs, the bumper—along with its attachment to the vehicle—absorbs the shock to prevent or reduce damage to the vehicle. 49 CFR Part 581, entitled “The bumper standard” and called simply “the standard” hereinafter, prescribes performance requirements for passenger cars in low-speed front (and rear) collisions. It applies to front (and rear) bumpers on passenger cars, and is aimed at preventing damage to the car body and safety related equipment at barrier impact speeds of 2½ mph across the full width and 1½ mph on the corners. This is equivalent to a 5 mph crash into a parked vehicle of the same weight. The standard requires protection in the region 16 to 20 inches above the road surface, and the manufacturer can provide the protection by any means it wants. For example, some vehicles do not have a solid bumper across the vehicle, but meet the standard by strategically placed bumper guards and corner guards. Some bumpers use energy absorbers or brackets and others are made with a foam cushioning material. Ordinary front (and rear) bumpers are not designed to be structural components that would significantly contribute to vehicle crashworthiness.\nAt this time, the U.S. military is involved in conflicts of a sort sometimes called urban warfare. The HMMWV (High Mobility Multi-purpose Wheeled Vehicle) is often used on patrol in such conflicts. As provided, the HMMWV is susceptible to being disabled by an ordinary passenger vehicle ramming into the front end of the HMMWV. The HMMWV as usually provided does not even have a front bumper. Also, the HMMWV is not usually armored, and when disabled, the crew is usually at a great disadvantage because they are usually outnumbered, surrounded, and without cover.\nIn addition, a strategy used in attacking a HMMWV on patrol is to place a wrecked or disabled junk vehicle in the roadway at a place where there is no way for the HMMWV to get around the junk vehicle. Crew members then would sometimes have to exit the HMMWV and push the junk vehicle out of the way, and would be shot at while trying to get the junk vehicle out of the way.\nWhat is therefore needed is a highly protective front bumper for the HMMWV and other vehicles used in such conflicts.\nThe prior art does provide highly protective front bumpers for such applications, but these are usually quite heavy, typically several hundred pounds, because of being made of steel. While there are advantages in having a steel front bumper, the HMMWV was not designed to maneuver with such a large amount of weight attached to the front.\nIn addition, the HMMWV is not a large vehicle, and is sometimes taken on extended duration patrols, and must then accommodate a large amount of supplies and equipment and several crew members. Storage space is therefore at a premium.\nThus, it would be advantageous to have a front bumper system that is not only high strength but also lightweight, and further, ideally, also includes some additional storage capacity."}
{"text": "This application relates generally to connectors and, more particularly, to connectors sized to receive a threaded apparatus.\nAircraft engines typically include a plurality of threaded connections that are torqued during assembly. For example, an aircraft engine may be mounted to an aircraft at various locations, such as on a wing, to a fuselage, or to a tail. The engine is typically mounted at a forward and an aft end with a mounting system that includes a plurality of thrust links including connectors at each end. The connectors include openings extending therethrough that are defined by substantially planar inner surfaces including a plurality of threads.\nThe thrust link connectors are exposed to potential vibrational stresses that are induced during engine operation. To minimize the effect of exposure to such vibrational stresses, threaded components are torqued. When threaded components are loaded by torque or an external load, stresses are induced within the connectors. A highest stress area typically occurs in a first thread. More specifically, a highest stress is induced into a first thread root fillet that is defined between the first thread and a subsequent thread. Continued exposure to such stresses, may weaken the first thread and reduce a low cycle fatigue capability of the threaded connection. Prolonged reduced low cycle fatigue may shorten a useful life of the threaded connection.\nTo facilitate improving the low cycle fatigue capability of threaded connections, the threads of at least some known threaded components are rolled to provide a surface that is pre-stressed in compression. Such a process is time-consuming. In some other known applications, thread rolling is not available and the threads are instead cut back to facilitate improving low cycle fatigue capability. However, such connections are also costly and time consuming to manufacture."}
{"text": "1. Field of the Invention\nThis invention relates to a breech closure for a barrel-type firearm, in which a flow passage for conducting an igniting jet from an igniting cartridge held in a magezine extends through a breechblock.\n2. Description of the Prior Art\nTo permit a gun to be fired after it has been loaded, an igniting cartridge must be inserted into the breech closure and a propellant charge for the projectile which has been loaded must be fired by an igniting jet, which is generated by the igniting cartridge and passes in a flow passage through the breechblock. Such igniting cartridges are usually inserted into the breech closure by hand because if a magazine is provided which is suitable for that purpose and in which the igniting cartridges are held it will be rather difficult to gastightly connect the magazine holding the igniting cartridges to the flow passage of the breechblock. In that case it must be taken into account that a gastightness in a very wide pressure range, e.g., between 100 and 4000 bars, must be ensured and that measures which may be adopted to seal against high pressures may sometimes be inconsistent with measures which will be required to seal against low pressures."}
{"text": "1. Field of the Invention\nThe present invention concerns a process for producing silicon carbide whiskers at high yield, at a reduced cost, with less content of powdery silicon carbide and at high purity.\n2. Description of the Prior Art\nSince silicon carbide whiskers are excellent in the strength, modulus of elasticity, oxidation resistance, heat resistance and chemical stability as compared with other fibers for use in composite materials, they have attracted attention as composite reinforcing material for ceramics, metals and plastics and, particularly, as composite reinforcing material for metals because they have satisfactory wettability with various kind of metals.\nThe production process for such silicon carbide whiskers is generally classified into (A) a method of depositing from liquid silicon carbide at high temperature and high pressure, (B) a method of dissolving carbon into a molten liquid of metal silicon and crystallizing silicon carbide, (C) a method of crystallizing by sublimating silicon carbide at high temperature from silicon carbide powder and (D) a method of growing crystals of silicon carbide formed by the heat decomposing reaction of silicon compounds.\nAmong them, the methods (A) and (B), however, result in a great difficulty in the production in view of the production facility because of the use of extremely high temperature and high pressure or the use of molten metal liquid. Further, the method (C) also requires an extremely high operation temperature and can not be operated with ease, as well as the facilities are complicated and fractional collection of the resultant whiskers is difficult. Accordingly, all of the methods (A) to (C) have fetal defects as the industrial production process.\nAccordingly, the method (D) as above in which silicon carbide is deposited in the form of whiskers while forming silicon carbide by various heat decomposing reactions is predominant at present as the process for producing silicon carbide and various methods have further been proposed therefor. That is, as the method (D), there have been known (1) a method of conducting solid phase reduction for silicon dioxide with carbon or metal silicon and carbon, (2) a method of reacting an organic silicon compound or a mixture of a silicon compound and a carbonaceous compound in a gaseous form at high temperature and (3) a method of gas phase growing utilizing the reaction between a fluorine-containing silicate and carbon.\nAmong them, the method (2) only produces extremely small amount of silicon carbide whiskers per unit volume of the reaction chamber since the reaction is conducted in a gas phase with a large volume of the chamber, requires to maintain the inside of the large capacity reaction chamber at a high temperature upon production in an industrial scale since the reaction is conducted at high temperature, resulting in high cost, as well as involves a problem that corrosive reaction by-products are formed. Furthermore, in the method (3), a fluorine-containing silicate is melted and reducing reaction with carbon added thereto is taken place and evolved vapors are cooled to form silicon carbide whiskers. Also the method (3) involves the same problems as that for the method (2) above in view of the utilization of gas phase reaction. Moreover, since the method (3) is necessary to handle a great amount of molten salt, it brings about problems that the material for the production facility is limited severely, and these molten salts are incorporated as impurity to the produced silicon carbide whiskers. Therefore, the method (3) is not preferable as the industrial production process for the silicon carbide whiskers.\nWhile on the other hand, in the method of utilizing the solid phase reduction (1) above, most easily available silicon dioxide is used as the starting material and it can be classified into two types depending on the reaction scheme.\n.circle.1 A process for producing silicon carbide whiskers by using silicon dioxide and carbon as the starting material by the following reaction: EQU SiO.sub.2 +C.fwdarw.SiO+CO EQU SiO+2C.fwdarw.SiC+CO\nThat is, fundamentally by the same method as the Acheson's method generally known as the industrial production process for silicon carbide.\n.circle.2 A production process by using silicon dioxide, metal silicon and carbon as the starting material by the reaction of: EQU SiO.sub.2 +Si.fwdarw.2SiO EQU SiO+2C.fwdarw.SiC+CO.sub.2\nAmong them, it is essential in the process .circle.2 to at first prepare silicon monoxide by the reaction of silicon dioxide and metal silicon. Accordingly, it is necessary, for example, as disclosed in Japanese Patent Publication No. Sho 47-18531 to preferentially conduct the reaction of forming silicon monoxide at first by silicon dioxide and metal silicon while avoiding the contact between silicon dioxide and carbon so that no unreacted metal silicon remains, for example, by a method of disposing a carbon or graphite substrate at a place apart from the portion of forming silicon monoxide. Accordingly, the process .circle.2 inevitably increases the capacity of the production facility and, also brings about defect of resulting in the reduction of the yield.\nWhile on the other hand, the process .circle.1 is free from the defect in the process .circle.2 described above and it is advantageous as the industrial production process since it can be conducted by merely mixing silicon dioxide and carbon uniformly and heating the mixture to a predetermined temperature.\nIn connection with the process .circle.1 , various studies have been made also to the starting silicon source material for producing silicon carbide whiskers by the process .circle.1 and there have been known, for example, (a) a method of using silicon component present in hulls of gramineous plants, (b) a method of using silica sand as the starting material, (c) a method of using silas or glass dust as the starting material, (d) a method of using active silica with high specific surface area such as silica sol and silica gel as the starting material, etc.\nHowever, the method (a) involves a drawback that hulls as the starting material have to be applied with pre-treatment of carbonizing or ashfying hulls and, in addition, a large-scaled hull processing facility is required since the silicon content of hulls is less than 20% by weight. Further, availability of hulls as the starting material is liable to suffer from the effects depending on the season and climate conditions and many impurities are incorporated into the silicon carbide whiskers produced. Further, the method (b) requires a step of previously pulverizing silica sands finely and the silica sands, being crystalline substance, show poor reactivity and results in lower yield as compared with amorphous silicic source starting material. The method (c) brings about problems that metal impurities other than silicon contained in the starting material may be incorporated into silicon carbide whiskers, or impurities are gasified and scattered in the heating step of the starting material, which are solidified by coagulation at the low temperature portion of the reaction chamber to hinder the continuous operation. The method (d) also involves problems that silica is expensive as the starting silicon source material, can provide no high reactivity as expected since the silica shrinks at high temperature and changes into deactivated silica with low specific surface area at the reaction temperature and, accordingly, provides no good yield. Therefore, it is difficult to provide silicon carbide whiskers with less impurity content at a high yield by using any one of starting silicon source materials (a) to (d).\nFurther, it has recently been disclosed a process for producing silicon carbide whiskers by using a hydrolyzable silicon compound as the starting silicon source material, admixing the hydrolyzable silicon compound and a carbon compound into a steam-containing hot gas, thereby forming a mixed aerosol of silicon oxide and carbon and producing silicon carbide whisker by using the mixed aerosol as the starting material (Japanese Patent Laid-Open No. Sho 60-16809 and Japanese Patent Publication No. Sho 61-2640). However, the process is disadvantageous in that the mixed aerosol formed by the process is extremely fine and of low bulk density to render the handling difficult, reduce the charging amount of the starting material per unit volume and worsen the productivity.\nFurthermore, in the process .circle.1 above, various attempts have been made for other materials than the starting silicon source material and the following improvement has also been proposed for the method of using the starting silicon source materials (b) to (d) described above.\n(i) A method of incorporating chlorine or hydrogen chloride, easily evaporizable hydrocarbon or halogenated hydrocarbon as a carbon source, etc. into an atmosphere for improving the reaction rate (Japanese Patent Publication No. Sho 52-28757 and U.S. Pat. No. 3933984).\n(ii) A method of adding, as VLS catalyst, powder of metals such as Fe, Co, Ni, B, La, Mn, Al, Ti, etc. or the compounds thereof to a mixture of silicon dioxide and carbon and, further, using vapors of sodium chloride as the space-forming agent (U.S. Pat. No. 3622272, Japanese Patent Publication Nos. Sho 51-8760, 59-45637 and 60-44280 and Japanese Patent Laid-Open No. Sho 61-22000).\nHowever, the method (i) involves a problems in using corrosive gases and the method (ii) involves problems in that it uses a great amount of the VLS catalyst and further uses sodium chloride, as the space-forming agent. Accordingly, the catalyst and sodium chloride are introduced as impurities to silicon carbide whiskers in the production step thereof to reduce the purity, thus bringing about the reduction in the wettability and adhesion between the resulting whiskers and the matrix, as well as reduction in the toughness of the whiskers.\nFurther, in any of the processes .circle.1 as described above, there is a high possibility that a great amount of powdery silicon carbide is by-produced together with the amied silicon carbide whiskers. Although the size and the content of the powdery silicon carbide vary depending on the whisker forming conditions, the powdery silicon carbide is by-produced while being broadly dispersed in the whiskers. When the whisker containing such by-products are composited with metals or plastics, the substantial whisker content composited is low and no expected compositing effect can be obtained. In addition, the strength of the composite material is rather reduced due to large grain size of granular or finely powderous product. Accordingly, the silicon carbide whiskers containing a great amount of powdery silicon carbide are often not suitable as composite reinforcing material for ceramics, metals and plastics. Furthermore, there has been a problem that the powdery products by-produced in the silicon carbide whiskers can not completely be separated even by the separation method utilizing the difference in the oleophilic and hydrophilic properties or by means of centrifugal precipitation.\nIn view of the foregoing situations, the present inventors have made earnest studies on the process for producing silicon carbide whiskers by improving the method (1) and, particularly, the process .circle.1 and have proposed a process of mixing a hydrolysis product of chlorosilane or chlorodisilane as the starting silicon source material with a carbonaceous material and reacting them at a high temperature from 1400.degree. to 1700.degree. C. (Japanese Patent Application Nos. Sho 60-26687 and Sho 61-53243). Although it is possible to produce silicon carbide whiskers at a yield as high as 80% or more, a further improvement is still required for suppressing the content of the by-produced powdery silicon carbide and reducing metal impurity."}
{"text": "Semiconductor memories are binary data memories in which a plurality of memory cells are provided. The memory cells are addressable by means of wordlines and bitlines. The main memory comprises a matrix of many memory cells connected to address decoding means and sense amplifiers. Addressing a memory cell, i.e. the selection of a memory cell, is performed by activating wordlines which are connected to address decoding means. The data stored in the addressed memory cells are read out by input/output units having sense amplifiers for amplifying the read-out data signal. The input/output units are connected to a data bus by means of which data can be read out from the memory or written into the memory.\nWhen a random access memory is produced, it can happen that some of the memory cells within the matrix are faulty. Accordingly, the produced memory chip is tested, and it is decided whether it can be repaired. For this purpose, there is normally provided an on-chip circuitry to provide testing of the memory chip. The built-in self-test (BIST) is essentially the implementation of logic built into the memory chip to do testing without the use of a tester for data pattern generation on comparison purposes.\nFIG. 1 shows the architecture of a memory chip according to the state of the art. The memory chip comprises a main memory with a plurality of memory cells and a redundancy logic having a built-in self-test device. The memory chip is connected to an address bus, a control bus and a data bus. In a test mode, the addresses of the faulty memory cells are detected. The redundancy logic replaces the faulty memory cells within the main memory with memory cells in the form of redundant registers within the redundancy logic. The address applied to the address bus is compared with the addresses of the detected faulty memory cells, and when there is a match, the redundancy logic maps the faulty address to an address of a register cell within the redundancy logic to replace the faulty memory cell. When reading data from the memory chip, the redundancy logic controls a multiplexer connected to the data bus. When accessing data with an address of a faulty memory cell, the data is read from the register replacing the memory cell within the redundancy logic.\nIn the conventional memory chip as shown in FIG. 1 the fusebox connected to the redundancy logic contains a plurality of fuses. The fuses are laser fuses which are blown by laser light. Alternatively a read only memory (ROM) is used instead of the fusebox. This ROM is for instance a EPROM or EEPROM. For blowing the fuses in the fusebox or for loading the data values into the read only memory the external tester device calculates which fuses have to be blown or what values have to be loaded into the ROM.\nTo repair a faulty RAM-memory within the memory chip several interactions are needed between the external tester device and the conventional BIST device as shown in FIG. 1. In a failure analysis the external tester device detects the addresses of the faulty memory cells within the RAM memory cell array and calculates a repair solution for repairing the memory cell array on the memory chip. The external tester device computes the fuses which have to be blown or respectively the values to be loaded into the ROM memory. After the repair solution has been calculated the fuses in the fusebox are blown by the external tester device.\nThe interactions between the external tester device and the memory chip are performed by using huge communication protocols. Consequently blowing of the fuses or loading of the data values into the read only memory takes a long time since the data exchange between the external tester device and the memory chip is extensive.\nWhen using laser fuses within fusebox for operation of the main memory a further disadvantage resides in that the laser fuses have to be blown by laser light. Accordingly to program the fuses it is necessary to have open access through the chip package to the laser fuses. In the testing process this leads to a complicated handling of the memory chips for reparation purposes.\nFIG. 2 shows the architecture of a main memory within the memory chip according to the state of the art. In this example, the main memory is a 8k×16 wide memory having 16 input/output units and wordline address decoders (XDEC) for decoding the wordline address or X-address of the memory cells. The input/output units are connected to the memory cell matrix by means of vertical bitlines. The input/output units receive the bitline address or Y-address of the solected memory cell.\nThe main memory shown in FIG. 2 according to the state of the art is partitioned in two memory hales wherein the X-address decoders are placed in the center. With this architecture, the length of the wordlines is comparatively short so that the parasitic capacitance of the wordlines can be minimized.\nEach input/output unit is connected to the 16 bitlines reading data from the addressed memory cell and for writing data into an addressed memory cell.\nFIG. 3 shows the architecture of an input/output unit according to the state of the art. For reading out data, the input/output unit comprises multiplexers which are connected to the bitlines of the memory cell matrix. In the shown example, each memory cell is connected to a multiplexer via a couple or pair of bitlines BL, {overscore (BL)} to provide a differential signal to the input of the multiplexer. In the shown example, each multiplexer has N signal inputs. On the output side, each multiplexer is connected to a differential amplifier and an inverter for amplifying the read-out data signal and to supply the data to a data bus. The multiplexers are controlled by the applied Y-address.\nIn a conventional memory, there are provided either redundant registers, redundant bitlines and/or wordlines to repair a memory chip in case that faulty memory cells are detected when testing the memory chip.\nIf the conventional memory chip comprises redundant registers, the number of faulty addresses is limited by the number of redundant registers provided within the redundancy logic. If there are, for instance, 10 redundant registers, it is only possible to repair 10 faulty addresses. When an address is “faulty”, the address is stored in a redundant register. Since the number of faulty addresses detected by the main memory, it is not known before testing a considerable number of registers have to be provided within the redundancy logic to guarantee the repair of the chip even when a lot of memory cells are detected to be faulty.\nIn case that the memory chip comprises redundant bitlines and/or wordlines, the repair method is much more complex, because all errors have to be known in advance before the error pattern can be diagnosed and an optimal repair solution can be calculated. Storing detected memories with a conventional method implies a very large array.\nSuch an array needs a lot of space on the memory chip, thus increasing costs when producing the memory chip.\nThe test diagnose repair performed with an external tester device as shown in FIG. 1 has the main disadvantages that it takes a long time to load the fuses within the fusebox because a very extensive data exchange is necessary between the external tester device and the memory chip. The communication between the external tester device and the memory chip is performed by using predetermined data protocols thus increasing redundancy of data exchange between the tester device and the memory chip. By increasing the reparation time of the memory chip the production costs of the memory chips are also increased.\nAccordingly it is the object of the present invention to provide a fuse blowing interface which allows to repair a memory chip within a minimum time span."}
{"text": "In treating hospitalized patients the body weight and its fluctuations is often of great importance, particularly in cases of actual or impending electrolyte and water imbalance. This determination supplements laboratory data, and makes its interpretation easier, more accurate, more relevant and more meaningful. It provides the clinical observer with objective information which may prove vital for proper and successful treatment of the patient.\nIn some cases, it is necessary that bedridden patients remain in their lying position. To weigh such patients the scale must be constructed so that the patient can be transferred from the bed onto the scale and weighed while he remains inclined. Such scales generally provide a support structure, usually freely movable over the ground by mounting the structure on casters and a weighing mechanism including a weigh-board carried for supporting the patient in his inclined position during the weighing operation. The weigh-board is sufficiently elevated above ground so that it can be moved over the patient's bed. The support structure includes a base disposed below the weigh-board which is moved beneath the bed while the weigh-board is positioned over the bed. To facilitate the storage of the scale while not in use the weigh-board is normally attached to the weighing mechanism so that it can be placed into an upright storage position. A scale of this type is disclosed in U.S. Pat. No. 3,032,131.\nThere are other prior art scales for weighing patients in their horizontal position which structurally differ from the one disclosed in the referenced patent. However, all scales have the general characteristic of providing a horizontal weigh-board, a movable support frame over the bed for receiving the patient while providing a sturdy and stable support for the weigh-board.\nPrior art scales for weighing patients in their horizontal position generally enable accurate measurements and they are relatively convenient for transferring the patient from the bed to the scale and back. However, they have a number of practical deficiencies which can render such scales cumbersome to handle and which may impair their accuracy unless the attendant carefully isolates the scale, and particularly the components of the weighing mechanism, from coming into contact with objects such as the side of the patient's bed, for example. If such contact does occur the measured weight is of course inaccurate. This is particularly serious because such inaccuracies cannot be detected after the measurement has been taken and the scales provide no means of warning the attendant that the weight he measures is inaccurate because the weighing mechanism is in contact with a foreign object, e.g., the bed.\nIn addition, prior art scales of the type under consideration have certain features which make them cumbersome to operate. For example, the weigh-board is normally pivotally mounted to a portion of the weighing mechanism so that it can be moved between an operative, horizontal position and a vertical storage position. To prevent the accidental movement of the board from one or the other position, it must be locked.\nComplicated locking arrangements for retaining the board in its upright storage position required careful operation of the locking mechanism on the part of the attendant. It was frequently necessary that the attendant walk around the scale to the point at which the lock can be applied. This is both cumbersome and time-consuming. An inattentive or careless attendant may forego the locking of the board altogether. While unlocked, however, the board may be accidentally tipped and pivoted into a horizontal position which is undesirable at least and which can injure bystanders or cause standing objects to be knocked down and damaged.\nAccordingly, prior art scales for weighing bedridden patients in their horizontal position, while they were adequate when properly operated, had several features which could compromise their accuracy and/or their safety."}
{"text": "There are a number of instances when it is desirable or necessary to test or monitor the concentration of an analyte, such as glucose, lactate, or oxygen, for example, in bodily fluid of a body. Bodily sample analyte tests are routinely conducted in a variety of medical settings (e.g., doctor's office, clinic, hospital, by medical personnel) and in the home by the patient and/or a caretaker. For example, it may be desirable to monitor high or low levels of glucose in blood or other bodily fluid that may be detrimental to a human. In a healthy human, the concentration of glucose in the blood is maintained between about 0.8 and about 1.2 mg/mL by a variety of hormones, such as insulin and glucagons, for example. If the blood glucose level is raised above its normal level, hyperglycemia develops and attendant symptoms may result. If the blood glucose concentration falls below its normal level, hypoglycemia develops and attendant symptoms, such as neurological and other symptoms, may result. Both hyperglycemia and hypoglycemia may result in death if untreated. Maintaining blood glucose at an appropriate concentration is thus a desirable or necessary part of treating a person who is physiologically unable to do so unaided, such as a person who is afflicted with diabetes mellitus.\nThe most important factor for reducing diabetes-associated complications is the maintenance of an appropriate level of glucose in the blood stream. The maintenance of the appropriate level of glucose in the blood stream may prevent and even reverse some of the effects of diabetes. Certain compounds may be administered to increase or decrease the concentration of blood glucose in a body. By way of example, insulin can be administered to a person in a variety of ways, such as through injection, for example, to decrease that person's blood glucose concentration. Further by way of example, glucose may be administered to a person in a variety of ways, such as directly, through injection or administration of an intravenous solution, for example, or indirectly, through ingestion of certain foods or drinks, for example, to increase that person's blood glucose level.\nRegardless of the type of adjustment used, it is typically desirable or necessary to determine a person's blood glucose concentration before making an appropriate adjustment. Typically, blood glucose concentration is monitored by a person or sometimes by a physician using an in vitro test that requires a blood sample. The blood sample may be obtained by withdrawing blood by lancing a portion of his or her skin, using a lancing device, for example, to make blood available external to the skin, to obtain the necessary sample volume for in vitro testing. The fresh blood sample is then applied to an in vitro sensor, such as an analyte test strip, which is positioned in a meter or the like, whereupon suitable detection methods, such as calorimetric, electrochemical, or photometric detection methods, for example, may be used to determine the person's actual blood glucose level.\nAvailable self-monitoring analyte systems generally include an analyte meter having a receptacle for receiving a cartridge containing a plurality of disposable analyte testing devices, e.g., test strips, which are individually and automatically dispensed on demand. An analyte test strip includes one or more chemical reagents designed to interact with a target analyte(s) in body fluid applied to it in such a way that an analyte meter connected to electrodes of the test strip can derive a value of the level of the target analyte contained in the body fluid. The meter may be an integrated device which provides a lancing function in addition to measuring the level of target analyte(s) in the body fluid being tested. Where the analyte test strips provide only a test component, the lancing is performed by a separate lancing mechanism housed within the meter. In other systems, the test strip may be integrated with a lancet into a single component, such as disclosed in U.S. patent application Ser. No. 12/488,181, herein incorporated by reference in its entirety. In either variation, the meter is configured to dispense the testing devices from the cartridge one at a time, as needed.\nThe test strip cartridge may include a biasing member at a loading end of the cartridge to bias the contained test strips towards the cartridge's strip-dispensing end. A seal is typically provided at the cartridge's strip-dispensing end for minimizing exposure of the test strips within the cartridge to ambient air. The seal is typically made of an elastomeric material configured to be released temporarily to permit loading of a test strip from the cartridge to within the meter for the lancing and testing processes. A desiccating material may be provided separately within the cartridge or integrated into the structural support of the cartridge to help maintain the analyte test devices substantially free of moisture. An advantage of an automatic test strip cartridge is that numerous analyte tests may be performed without having to manually load a new test strip for each test performed.\nExamples of integrated lancing and testing systems are described in U.S. patent application Ser. Nos. 10/629,348, 10/701,993, 10/837,886, 10/899,773, 11/035,131, 11/146,897, 11/160,407, 11/160,427, 11/350,398, 11/535,985, 11/535,986, 11/830,760, 11/830,770, 11/830,779, 11/830,786, 11/831,649, 11/831,706, 11/868,762, 11/870,420 and 12/035,348, the disclosures of which are herein incorporated by reference.\nDue to certain lot-to-lot inconsistencies in the test strip fabrication and manufacturing process, there may be variations in test strip sensitivity between lots which require some form of system calibration for each batch of strips. Currently, existing calibration mechanisms require the use of a calibration strip by the user, the inputting of a calibration code by the user, or the use of a machine readable mechanism on the strip or cartridge to modify the reaction interpretation of the meters for a particular lot of test strips.\nDespite means for calibrating meters for given test strip lots, if a test strip is past its expiration date or its chemistry has otherwise degraded due to humidity and other constituents of the ambient atmosphere to which it has been exposed subsequent to fabrication and packaging and prior to actual use, i.e., during the test strip's shelf-life, unpredictable and unreliable analyte test results are likely. Moreover, during actual active use of a test strip cartridge, the cartridge is opened and closed by the meter each time a test strip is used, exposing the remaining test strips to an even higher level of ambient air, thereby contributing to degradation of the strips, even in the presence of a desiccating material within the cartridge."}
{"text": "In power plants a surface condenser has been conventionally employed to cool and condense a gas (steam) exhausted from a power generating steam turbine for recycling the condensed water. In a condenser of this type thousands of or tens of thousands of copper alloy condenser tubes, having ordinarily an internal diameter of 10-40 mm.phi. and a length as long as 5-40 m, are incorporated. Through the condenser tubes flows cooling water such as sea water so as to cool down the exhaust steam (vapor) passing the outside gap of the large number of condenser tubes.\nThe copper alloy condenser tubes are subjected to various types of corrosion because of fairly high speed flowing of corrosive water such as sea water on the order of 1-2.5 m/s and of the pollution of cooling water, for example, which necessitates coating or painting of the inner surface of the tubes with an anti-corrosive paint of organic resin over the whole length thereof for the purpose of preventing corrosion. However, the coating is strictly conditioned not to greatly inhibit the heat transfer rate of copper alloy tubes. In the case of applying anti-corrosion coating on the inner surface of copper alloy condenser tubes, the thickness of the coated film must be uniform and controlled in the range as small as 10-30.mu. from the view point of maintaining the desired heat transfer rate.\nThe life of coating of this type is inevitably shorter than the life of the power plant itself, which is as long as 20-30 years if the above-mentioned thin coated film is maintained; and the coating is liable to be worn away after the tube is mounted in the plant to expose sometimes the base metal. The coating may be damaged in some cases by mechanical rubbing with shells or sand grains contained in the cooling water, and its wearing may be accelerated by the so-called sponge-ball cleaning carried out to remove mud and/or sea weeds stuck on the inner surface of the tubes. Such being the case, the tubes must be recoated periodically or occassionally to maintain anti-corrosion and anti-rust condition in the plant.\nSome of coating methods have been recognized and practiced widely for coating the interior of tubes of relatively small length, for example, flowing paint in a tube or directly brushing paint. Such methods are however utterly impracticable, in the case of coating a long tube of small diameter such as a condenser tube, for obtaining a coated film of uniform thickness there. And particularly in the former method the paint flowed into a tube can not be diffused in the interior of the tube unless the horizontal tubes installed in the condenser are inclined.\nAs a relatively practical method spray coating can be mentioned, wherein the interior of a tube is coated by a spray gun spraying atomized paint. Even in this method employing a long necked spray gun of 500 mm, a coatable area, or the length of the coated area, covered by means of inserting the gun into the tube is naturally limited (restricted) in the partial length thereof. As a variation of such spray coating, a method of coating, wherein a paint spraying nozzle is moved (retracted) from one end opening of a tube to the other end opening during spraying paint, has been developed and has succeeded in getting a coated film of uniform thickness throughout the whole length of a long tube of small diameter.\nThere have still been, however, several problems in the spray coating method of this type: it is a matter of course that no spray is allowed while the spray gun is moved through a long tube from one end to the other end before it reaches the destination where spraying is to begin; even a slight leakage of paint from the nozzle in the meantime will hamper the uniform coating of the tube interior; the paint passage must be absolutely tightly closed while it is not in use, otherwise the nozzle will suffer from uneven spraying or no-spraying due to hardening of paint left in the nozzle.\nFor the elimination of those problems conventional nozzles have employed a stopper of core bar type. In other words, a pointed portion on the tip of the stopper is fitted into an opening portion of the paint passage of a nozzle insert, so that the opening portion of the paint passage is opened and closed by the advancing and retracting of the stopper. Such a type of stopper is liable to abnormally work or in some cases become entirely paralysed, in the event that a tube to be coated is very long, due to unexpected deflection or friction of the stopper.\nThe inventors of this invention proposed, in JITSU-GAN-SHO-54(1979)-147332 and JITSU-GAN-SHO-54(1979)-147333 (Japan, etc., a method of eliminating the above-mentioned disadvantages by means of closely and tightly covering the paint spraying nozzle by a stopper of a cap shape type. In a nozzle of this type the tip portion thereof is covered by a cap type stopper until it reaches the other end of a tube to be coated. Upon the nozzle reaching the other end of the tube, the cap is removed from the nozzle to allow the paint to be sprayed while the nozzle is being retracted along through the tube inside. This enabled coating of the interior of a long tube. It still leaves, however, something to be desired, such as the necessity of arranging an individual operator on either end of the tube to be coated for the purpose of putting on and removing the cap type stopper, and particularly in the event of coating copper alloy tubes in a condenser at a power station the stopper removed on one side of a condenser must be brought back each time to the other side of a condenser where the nozzle is inserted. It not only provides a great problem of coating efficiency, but also another serious problem of deteriorating the operation environment owing to a remarkable rising of density of the harmful organic solvent, such as toluene, on the stopper removing side at the beginning of paint spraying.\nIn the spray coating of copper alloy tubes in a condenser there are some unavoidable restriction conditions, especially in the event of coating or re-coating of the already installed tubes in place, from the view point of operation mode, operation environment, or operation term allowed, etc. A first problem of restriction lies in an extremely narrowly limited space for the coating operation; it is often limited in the condenser water box, the dimensions of which are 2-3 m in depth, 2-3 m in width, and 2-5 m in height. In the case of coating tubes in such a restricted space, ordinary operational apparatus, techniques and conditions used in a satisfactorily broad manufacturing plant can not be applied as they are.\nA second problem of restriction lies in deterioration of the operational environment, due to gradual pollution of the atmosphere in the condenser water box where the coating operation must be carried out, because the evaporation of thinner (solvent) from the paint in the narrow operation space makes it undesirable to keep the operators staying there for a long time.\nA third problem of restriction lies in that the re-coating operation must be finished during the term of inspection of the whole power generating plant, so the term allowed for the operation of re-coating the heat exchanger tubes is relatively short. Even when an established operation mode in an already installed plant may be introduced there, the number of coating apparatuses permitted to be worked in the narrow space is naturally limited. Elongation of the operation term is very inconvenient under such circumstances. A coating method and apparatus of high efficiency has thus been badly looked forward to."}
{"text": "This invention provides isolated DNA sequences, proteins encoded thereby, and methods of using said DNA and protein in a variety of applications.\nWidespread antibiotic resistance in common pathogenic bacterial species has justifiably alarmed the medical and research communities. Frequently, resistant organisms are co-resistant to several antibacterial agents. Penicillin resistance in Streptococcus pneumoniae has been particularly problematic. This organism causes upper respiratory tract infections. Modification of a penicillin-binding protein (PBP) underlies resistance to penicillin in the majority of cases. Combating resistance to antibiotic agents will require research into the molecular biology of pathogenic organisms. The goal of such research will be to identify new antibacterial agents.\nWhile researchers continue to develop antibiotics effective against a number of microorganisms, Streptococcus pneumoniae has been more refractory. In part, this is because Streptococcus pneumoniae is highly recombinogenic and readily takes up exogenous DNA from its surroundings. Thus, there is a need for new antibacterial compounds and new targets for antibacterial therapy in Streptococcus pneumoniae."}
{"text": "It has become common for users of host computers connected to the World Wide Web (the \"Web\") to employ Web browsers and search engines to locate Web pages having specific content of interest to users. A search engine, such as Digital Equipment Corporation's AltaVista search engine, indexes hundreds of millions of Web pages maintained by computers all over the world. The users of the hosts compose queries, and the search engine identifies pages that match the queries, e.g., pages that include key words of the queries. These pages are known as a result set.\nIn many cases, particularly when a query is short or not well defined, the result set can be quite large, for example, thousands of pages. The pages in the result set may or may not satisfy the user's actual information needs. Therefore, techniques have been developed to identify a smaller set of related pages.\nIn one prior art technique used by the Excite search engine, please see \"http://www.excite.com,\" users first form a query that attempts to specify a topic of interest. After the result set has been returned, the user can use a \"Find Similar\" option to locate related pages. However, there the finding of the related pages is not fully automatic because the user first is required to form a query, before related pages can be identified. In addition, this technique only works on the Excite search engine and for the specific subset of Web pages that are indexed by the Excite search engine.\nIn another prior art technique, an algorithm for connectivity analysis of a neighborhood graph (n-graph) is described by Kleinberg in \"Authoritative Sources in a Hyperlinked Environment,\" Proc. 9th ACM-SIAM Symposium on Discrete Algorithms, 1998, and also in IBM Research Report RJ 10076, May 1997, see, \"http://www.cs .cornell.edu/Info/People/kleinber/auth.ps.\" The algorithm analyzes the link structure, or connectivity of Web pages \"in the vicinity\" of the result set to suggest useful pages in the context of the search that was performed.\nThe vicinity of a Web page is defined by the hyperlinks that connect the page to others. A Web page can point to other pages, and the page can be pointed to by other pages. Close pages are directly linked, farther pages are indirectly linked. This connectivity can be expressed as a graph where nodes represent the pages, and the directed edges represent the links. The vicinity of all the pages in the result set is called the neighborhood graph.\nSpecifically, the Kleinberg algorithm attempts to identify \"hub\" and \"authority\" pages in the neighborhood graph for a user query. Hubs and authorities exhibit a mutually reinforcing relationship.\nIn U.S. patent application Ser. No. 09/007,635 \"Method for Ranking Pages Using Connectivity and Content Analysis\" filed by Bharat et al. on Jan. 15, 1998, a method is described that examines both the connectivity and the content of pages to identify useful pages. However, the method is relatively slow because all pages in the neighborhood graph are fetched in order to determine their relevance to the query topic. This is necessary to reduce the effect of non-relevant pages in the subsequent connectivity analysis phase.\nIn U.S. patent application Ser. No. 09/058,577 \"Method for Ranking Documents in a Hyperlinked Environment using Connectivity and Selective Content Analysis\" filed by Bharat et al. on Apr. 9, 1998, a method is described which performs content analysis only a small subset of the pages in the neighborhood graph to determine relevance weights, and pages with low relevance weights are pruned from the graph. Then, the pruned graphed is ranked according to a connectivity analysis. This method still requires the result set of a query to form a query topic.\nIn any of the above cases, it would be advantageous if duplicate or near duplicate pages could quickly be identified since these pages essentially represent the same content. It would even be better, if near duplicates could be identified without having the analyze the detailed content of the pages."}
{"text": "Radio communications receivers commonly receive and decode messages which contain information. A part of this information may be an identity for the intended recipient or recipients of a given call. A radio, upon decoding valid identity information, will receive a call. For example a radio may unsquelch its receiver and pass audio signals to a loudspeaker upon receipt of a valid identity. Conversely, radios commonly transmit a message having identity information so that other radios intended to receive a given call can receive and decode the transmitted identity information.\nModern radios transmit and receive messages in the form of digital data. Such dale is modulated onto a radio frequency carrier which is used to accomplish radio communications. A portion of any given message may he the aforementioned identity information. This information is commonly called an ID code. A transmit ID code is transmitted by the transmitting radio, which is received as a receive ID code by the receiving radio. A single given radio may or may not use different transmit ID codes and received ID codes, depending on what other radios the given radio is attempting communications with.\nA radio that is capable of both transmitting and receiving is known as a two-way radio. Such radios may operate in the simplex, half-duplex, and full-duplex modes of operation. These modes of operation are well known in the art. A simplex radio transmits and receives on one frequency during a call and is not capable of simultaneously transmitting and receiving. Rather, transmissions and receptions must be multiplexed in time in order to have a conversation.\nA half-duplex radio is capable of transmitting and receiving on a pair of two different frequencies, and, like a simplex radio is not capable of simultaneously transmitting and receiving. A full-duplex radio operates on a pair of two different frequencies and is capable of simultaneously transmitting and receiving.\nBoth half-duplex and full-duplex radios are capable of operating through a radio repeater. A repeater is a radio transmitter and receiver pair which receive signals from radios and retransmit them to other radios. Radio repeaters, or repeaters, typically include a control unit for controlling their operation. One benefit of a repeater is that it can be located at a high elevation and operated at high power. This serves to extend the radio coverage range to a greater geographic area than would be available if two radios were to communicate directly with each other. For example, radios are typically deployed in the form of mobile radios for use in motor vehicles and the like or in the form of hand-held portable radios. In either case, size, power and elevation limitations would prevent these radios from communicating over long ranges or large geographic areas.\nOften times, repeaters are shared between a large number of radio users. Many times, radios users are grouped together so that users with common communications needs can effectively communicate with each other. In some situations two radios are grouped together to provide for one to one communications. Alternatively, several radios may be grouped together. In situations where several radios are grouped together, a single radio will transmit messages to the remainder of the group. These radio groupings are called radio fleets. The one to one calls are called private calls and the one to many calls are called dispatch calls. A single radio may be a member of more than one group. For example, a single radio user may desire to have private calls in some instances and dispatch calls in other instances. Furthermore, any given radio may be a member of several private groups and/or dispatch groups. Another application of radios and repeaters is to interconnect a repeater to the Public Switched Telephone Network, (PSTN). With an interconnected repeater, a radio may communicate with the repeater via radio signals and further communicate with a telephone user via wires in the PSTN. Radios operating in the interconnected environment can place call into the PSTN or receive calls that originate in the PSTN. These calls are commonly known as interconnect calls.\nIn radio environments where there are a very large number of individual radio users, a single repeater may not be sufficient to provide an adequate degree of service. The Federal Communications Commission rules allow for the trunking of two or more repeaters together to form a trunked groups of repeaters, also called a repeater system. These repeater systems employ a trunking signaling protocol which allow the management of individual repeaters in real time. This trunking information is often times a part of the radio messages transmitted between radios and repeaters.\nTransmit ID codes and receive ID codes are used to identify the transmitting and receiving radios in the various groupings as described above. Transmit and receive ID codes are usually transmitted as binary numbers and they are sometimes represented as base ten numbers. System operators who program ID codes, sell radios, and service radios become technically experienced with ID code usage and their limitations. However, to end users of radios, these ID codes can be confusing and difficult to work with. Yet it is still important for the end user to identify which dispatch fleet or private group the end user is calling or is to be called. Radios usually have a selection means, such as a push button or rotary knob for selecting a communications group which is specified by a group ID code. An association is often made between group ID codes and transmit ID codes and receive ID codes so that end users can simply select a familiar, simple group ID code which is ultimately linked to the appropriate transmit ID and receive ID needed to accomplish the desired communications.\nIn one example, a simple system has been devised that links group ID codes one through ten to pairs of a transmit ID's and a receive ID's. The radio user simply remembers, for example, that group 1 places a certain private call, group 2 places a telephone interconnect call, group 3 places a dispatch call, and etc. While receiving calls, the users radio may be enabled to scan all of the aforementioned receive ID's so that receipt of any one of them in a data message will cause the call to be received. When a call is received, the radio may display the familiar group ID code so that the radio user can identify the calling group. 0f course, other transmit ID codes and receive ID codes could be used, such as alpha-numeric characters, for example, that identify radio groups with words, numbers or mnemonics.\nIn partial summary, group ID codes are used to specify and identify specific transmit ID's and receive ID's for the purpose of controlling what group or fleet is being communicated with.\nA described earlier, repeaters offer radio coverage over a larger geographic area than is possible with direct radio to radio communications. However, this does not mean that the coverage area is without limits. For various reasons that are well know in the art, repeaters also have finite coverage area. To offer even larger coverage area, operators of radio repeater systems may operate two or more repeater systems and offer service to individual radios on more that one repeater system.\nMultiple repeater systems are often located in separate geographic areas. As a user of a radio moves from area to area, it is necessary to select the repeater system on which to initiate calls. Conversely, the radio must be tuned to a frequency of one of the repeaters in the local system in order to receive a call. Radios have been devised with a system scan feature which causes the radio to scan several frequencies of repeaters that are located in different repeater systems. While monitoring any particular repeater frequency, the radio also scans for particular receive ID codes which would indicate that a desired call is to be received.\nIn a similar fashion to the identification of particular transmit and receive ID codes with a group ID code, repeater systems are identified with a simple system ID code. For example, system ID codes may be numbered from one through ten to identify up to ten separate repeater systems, each of which may have one or more repeaters. Again, the radio user simply need remember, for example, that system 1 is used downtown, system 2 is used north of town, system 3 east of town, and etc.\nIt is now understood that a user of a radio selects a system to identify the repeater system, or geographic area, to communicate in, and a group to identify with whom the communications is to be had. Likewise when a call is received, a radio may display the system and group numbers which will identify the repeater system being used and which particular radio or group of radios is calling.\nModern radios have controllers which allow them to function with very little user intervention. Calls can be received automatically and the response to a call usually only requires the press of a single button. Such an automated system functions well if the radio user is present at the radio to communicate in a conversation. However, often times, the radio user is away from the radio while a call is received. Even while the user is away from the radio, the radio will scan systems and groups, received and decode data messages into ID codes, unsquelch the receiver, receive the call, and etc. However, all this information is meaningless if the users does not listen to, interpret, and/or respond to it. Upon the user's return to the radio, the user is unaware of the missed call, and therefore, communications is not completed. Since the main benefit of two-way radios is to improve the accessibility to communications for the end user, and since a missed called reduces the desired accessibility, clearly there is a need to provide a greater amount of call history detail to users of two-way radios, thus improving the accessibility to communications for the end user."}
{"text": "The present invention relates to a laser beam machining apparatus used in processing a printed circuit board or the like.\nConventionally available laser beam machining apparatuses are based on either a method in which processing is effected at a focal position of a laser beam by merely focusing the beam or a method in which processing is effected by inserting a mask into a laser beam transmission path and by reducing an image of the mask and forming it on a processing surface after setting the beam diameter to a desired size.\nFIG. 27 is a diagram illustrating a cross section of a multilayered printed circuit board. The multilayered printed circuit board is a printed circuit board in which layers of a printed circuit board are insulated from each other by a resin and are superposed one on top of another. To interconnect wiring portions (copper foils) of one layer and another layer, blind via holes or through holes are formed in the substrate, and the cross sections of the holes are plated, thereby interconnecting the wiring portions between the layers. The diameters of the blind via holes or the through holes are becoming smaller due to the trend toward a higher degree of the packaging density of electronic components. In recent years, hole diameters on the order of several hundred microns have come to be required, and the processing of holes having different diameters have been required for one substrate. With the hole diameters of such a measure, it is difficult to cope with them by the conventional drilling, so that processing by using a laser beam is required. In addition, as the increasing trend toward a greater number of layers in the printed circuit board, the processing of blind via holes having different hole depths is required for one substrate.\nFurther, due to the higher degree of the density of the multilayered printed circuit board, cut processing has also come to be effected by a laser beam.\nTo effect the fine hole processing by using a laser beam, it is necessary to reduce the size of the laser beam to the size of the diameter of the hole to be processed. As such methods, the following two methods are known: a method in which the laser beam is focused by a lens (the system based on this method will be hereafter referred to as a focusing optical system) and a method in which a mask is provided in the optical path of the laser beam (laser beam transmission path) and an image of the mask is reduced and formed on a processing surface by means of a lens (the system based on this method will be hereafter referred to as an image-transfer optical system).\nNext, a description will be given of the focusing optical system and the image-transfer optical system.\nFIG. 28 is a diagram illustrating a conventional focusing optical system. In the focusing optical system, a laser beam is focused by a lens, and processing is effected in the vicinity of the focal position of the laser beam. At that time, a spatial filter such as the one shown in FIG. 29 is sometimes used to eliminate aberration components of the laser beam for the purpose of improving the laser-beam focusing characteristic.\nIn FIG. 28, reference numeral 1 denotes a laser oscillator; 2, a laser beam; 5, a focusing optical system; 6, an XY table; 7, a workpiece; 8, a reflecting mirror; and 9, a numerical controller (hereafter referred to as the NC device). The focusing optical system 5 is shown in FIG. 29. Numeral 51 denotes a lens; 52, a spatial filter; and 53, a focusing lens. The spatial filter 52 is disposed at a focal position of the lens 51, i.e., at the position where the beam is subjected to Fourier transform. The spatial filter 52 is used to enhance the beam focusing performance by eliminating the aberration components of the beam, and a detailed description is given in, for example, Walter Koechner, \"Solid-State Laser Engineering,\" Springer-Verlag, 1992, pp. 174-180.\nNext, a description will be given of the operation.\nA laser beam emitted from the laser oscillator 1 is subjected to Fourier transform onto the spatial filter 52 by the lens 51, and the spatial filter 52 transmits only low-frequency components of the spatial frequencies of the laser beam. The beam with its aberration components eliminated by the spatial filter 52 is focused on the surface of the workpiece 7 by the focusing lens 53.\nA plurality of processing conditions, i.e., such conditions as a laser output and the like, are stored in advance in a storage device of the NC device 9. Optimum conditions suited to the material, the thickness; and the processing shape of the workpiece 7 are selected, and the oscillator 1, the XY table 6, and the like are controlled on the basis of the selected conditions.\nFIG. 30 is a diagram illustrating a conventional image-transfer optical system. A mask is provided in the laser beam transmission path, and an image of a pin hole in the mask is reduced and formed on a processing surface by means of a lens so as to obtain at the processing surface the laser beam having a diameter defined by the mask.\nIn FIG. 30, reference numeral 1 denotes the laser oscillator; 2, the laser beam; 3, an image-transfer optical system; 6, the XY table; 7, the workpiece; 8, the reflecting mirror; and 9, the NC device. The image-transfer optical system 3 is shown in FIG. 31. Reference numeral 31 denotes a lens; 32, a mask; and 33, a lens. The positional relationship among the mask 32, the lens 33, and the workpiece 7 is in such a positional relationship that the image of the mask is formed on the workpiece 7 at a certain magnification. As for image formation, a detailed description is given in, for example, K. Iizuka, \"Engineering Optics,\" Springer-Verlag, 1985, pp. 145-164.\nNext, a description will be given of the operation.\nA laser beam emitted from the laser oscillator 1 is made incident upon the image-transfer optical system 3, and is focused by the lens 31 so as to be applied to the mask 32. The reason for focusing is to reduce the energy loss in irradiating the mask having a small diameter. As for the beam which is transmitted through the mask 32, an image of the mask is formed on the surface of the workpiece 7 at a certain reduction ratio by the lens 33.\nA plurality of processing conditions, i.e., such conditions as a laser output and the like, are stored in advance in the storage device of the NC device 9. Optimum conditions suited to the material, the thickness, and the processing shape of the workpiece 7 are selected, and the oscillator 1, the XY table 6, and the like are controlled on the basis of the selected conditions.\nIn the image-transfer optical system, to change the beam diameter on the workpiece 7, a plurality of different masks are prepared in advance, and the mask is replaced by a different one when the beam diameter on the workpiece 7 is to be changed.\nImportant parameters in the fine-hole processing by using a laser beam are a minimum critical hole diameter R and a maximum critical hole depth DOF. The minimum critical hole diameter R and the maximum critical hole depth DOF can be expressed by the following formulae denoted by the singular reference number (1): EQU R=k1.times.l.times.F EQU DOF=k2.times.l.times.(F.sup.2) (1)\nwhere\nF=D/f PA1 k1 (focusing optical system)&gt;k1 (image-transfer optical system) PA1 k2 (focusing optical system)&gt;k2 (image-transfer optical system) PA1 R (focusing optical system)&gt;R (image-transfer optical system) PA1 DOF (focusing optical system)&gt;DOF (image-transfer optical system)\nD is an effective diameter of the lens, f is a focal length of the lens, l is a wavelength of the laser beam, k1 and k2 are values which are determined by the material of the workpiece 7 and the state of the laser beam, i.e., the amount of aberration and the mode. A comparison of the values of k1 and k2 between the focusing optical system and the image-transfer optical system is shown below. However, it is assumed that the material of the workpiece 7 is the same, and that the state of the beam is also the same.\nThat is,\nIn processing, the image-transfer optical system is capable to processing a hole having a smaller diameter than the focusing optical system, but the depth of the hole which can be processed is shallower in the case of the image-transfer optical system than in the case of the focusing optical system.\nFIG. 32 shows processable regions in the focusing optical system and the image-transfer optical system with respect to a certain material subjected to hole processing.\nAccordingly, the processing of holes of small diameters is conventionally effected by a laser beam machining apparatus provided with an image-transfer optical system, while the processing of holes of large depths is conventionally effected by a laser beam machining apparatus provided with a focusing optical system. In addition, in the cut processing of multilayered printed circuit boards, since a beam having a large focal depth rather than a beam having a small diameter at the processing surface is required, cut processing is conventionally effected by a laser beam machining apparatus provided with the focusing optical system.\nFor this reason, in the processing of one multilayered printed circuit board, in order to efficiently and consistently effect the cut processing as well as the processing of through holes and blind via holes having different diameters and depths, it is essential to change over the system between the focusing optical system and the image-transfer optical system in one laser beam machining apparatus.\nIn the processing of a multilayered printed circuit board using a laser beam, to enhance productivity, an attempt is made to realize high-speed production by moving not only the XY table but also the beam by using a galvanomirror. FIG. 33 shows a laser beam machining apparatus using a galvanomirror.\nIn FIG. 33, reference numeral 1 denotes the laser oscillator; 2, the laser beam; 19, an image-transfer optical system; 6, the XY table; 7, the workpiece; 8, the reflecting mirror; and 9, the NC device. The image-transfer optical system 19 is shown in FIG. 34. Reference numeral 31 denotes the lens; 32, the mask; 36, an f.theta. lens; 37, a galvanomirror. The positional relationship among the mask 32, the lens 36, and the workpiece 7 is in such a positional relationship that the image of the mask is formed on the workpiece 7 at a certain magnification when the beam is directed vertically downward from the galvanomirror 37 to the f.theta. lens 36.\nNext, a description will be given of the operation.\nA laser beam emitted from the laser oscillator 1 is made incident upon the image-transfer optical system 19. The beam made incident upon the image-transfer optical system 19 is focused by the lens 31 so as to be applied to the mask 32. The reason for focusing is to reduce the energy loss in irradiating the mask having a small diameter. The beam which has been transmitted through the mask 32 is led to the transferring f.theta. lens 36 by the galvanomirror 37, and an image of the mask reduced at a certain magnification by the f.theta. lens 36 is formed on the surface of the workpiece 7. The galvanomirror 37 is capable of guiding the beam to an arbitrary position on the f.theta. lens 36, with the result that the beam can be scanned by the galvanomirror 37 to the extent corresponding to the region of the size of the f.theta. lens 36 on the workpiece 7.\nA plurality of processing conditions, i.e., such conditions as a laser output and the like, are stored in advance in the storage device of the NC device 9. Optimum conditions suited to the material, the thickness, and the processing shape of the workpiece 7 are selected, and the oscillator 1, the XY table 6, the galvanomirror 37, and the like are controlled on the basis of the selected conditions.\nThe output of a laser oscillator is stable in the vicinity of a rated output, but at a low output the fluctuation width becomes large and the output becomes unstable. Since only several pulses are used in hole processing, variations in the laser output exerts a large effect on variations in the result of processing. Accordingly, the variations in the result of processing become large in a low-output region where the fluctuation width is large.\nSince the conventional laser beam machining apparatus is constructed as described above, a number of problems that are listed below have been encountered.\nThe laser beam machining apparatus having only the focusing optical system is incapable of processing small-diameter holes since its minimum critical processing hole diameter is larger than that in the case of the laser beam machining apparatus having only the image-transfer optical system. Meanwhile, the laser beam machining apparatus having only the image-transfer optical system is incapable of processing deep through holes and blind via holes since the depth of the holes which the apparatus can process is shallower than that in the case of the laser beam machining apparatus having only the focusing optical system. In addition, in cut processing, if the thickness of the printed circuit board becomes large, a beam having a large focal depth at the processing surface is required, so that processing cannot be effected by the laser beam machining apparatus having only the image-transfer optical system.\nFurther, in the image-transfer optical system, the mask 31 must be changed in order to change the beam diameter on the workpiece 7. Accordingly, productivity declines by the portion of the time need for mask replacement. Additionally, even if a number of kinds of masks with different hole diameters are prepared, it is impossible to cope with the processing of holes with continuously changing diameters.\nIf a comparison is made between a case where the beam is directed vertically downward from the galvanomirror 37 to the f.theta. lens 36 and a case where it is directed diagonally downward, the distance between the mask 32 and the f.theta. lens 36 differs by L1-L0, as shown in FIG. 35. Consequently, in the case where the beam is directed diagonally downward from the galvanomirror 37 to the f.theta. lens 36, the image of the mask 32 on the workpiece 7 becomes blurred, so that satisfactory processing cannot be effected.\nIn the low-output region where the fluctuation width is large, variations in the result of processing become large."}
{"text": "The detection and characterization of specific nucleic acid sequences and sequence changes have been utilized to detect the presence of viral or bacterial nucleic acid sequences indicative of an infection, the presence of variants or alleles of mammalian genes associated with disease and cancers, and the identification of the source of nucleic acids found in forensic samples, as well as in paternity determinations.\nVarious methods are known in the art which may be used to detect and characterize specific nucleic acid sequences and sequence changes. Nonetheless, as nucleic acid sequence data of the human genome, as well as the genomes of pathogenic organisms accumulates, the demand for fast, reliable, cost-effective and user-friendly tests for specific sequences continues to grow. Importantly, these tests must be able to create a detectable signal from a very low copy number of the sequence of interest. The following discussion examines three levels of nucleic acid detection currently in use: I. Signal Amplification Technology for detection of rare sequences; II. Direct Detection Technology for detection of higher copy number sequences; and III. Detection of Unknown Sequence Changes for rapid screening of sequence changes anywhere within a defined DNA fragment.\nI. Signal Amplification Technology Methods for Amplification\nThe “Polymerase Chain Reaction” (PCR) comprises the first generation of methods for nucleic acid amplification. However, several other methods have been developed that employ the same basis of specificity, but create signal by different amplification mechanisms. These methods include the “Ligase Chain Reaction” (LCR), “Self-Sustained Synthetic Reaction” (3SR/NASBA), and “Qβ-Replicase” (Qβ).\nPolymerase Chain Reaction (PCR)\nThe polymerase chain reaction (PCR), as described in U.S. Pat. Nos. 4,683,195 and 4,683,202 to Mullis and Mullis et al., describe a method for increasing the concentration of a segment of target sequence in a mixture of genomic DNA without cloning or purification. This technology provides one approach to the problems of low target sequence concentration. PCR can be used to directly increase the concentration of the target to an easily detectable level. This process for amplifying the target sequence involves introducing a molar excess of two oligonucleotide primers which are complementary to their respective strands of the double-stranded target sequence to the DNA mixture containing the desired target sequence. The mixture is denatured and then allowed to hybridize. Following hybridization, the primers are extended with polymerase so as to form complementary strands. The steps of denaturation, hybridization, and polymerase extension can be repeated as often as needed, in order to obtain relatively high concentrations of a segment of the desired target sequence.\nThe length of the segment of the desired target sequence is determined by the relative positions of the primers with respect to each other, and, therefore, this length is a controllable parameter. Because the desired segments of the target sequence become the dominant sequences (in terms of concentration) in the mixture, they are said to be “PCR-amplified.”\nLigase Chain Reaction (LCR or LAR)\nThe ligase chain reaction (LCR; sometimes referred to as “Ligase Amplification Reaction” (LAR) described by Barany, Proc. Natl. Acad. Sci., 88:189 (1991); Barany, PCR Methods and Applic., 1:5 (1991); and Wu and Wallace, Genomics 4:560 (1989) has developed into a well-recognized alternative method for amplifying nucleic acids. In LCR, four oligonucleotides, two adjacent oligonucleotides which uniquely hybridize to one strand of target DNA, and a complementary set of adjacent oligonucleotides, which hybridize to the opposite strand are mixed and DNA ligase is added to the mixture. Provided that there is complete complementarity at the junction, ligase will covalently link each set of hybridized molecules. Importantly, in LCR, two probes are ligated together only when they base-pair with sequences in the target sample, without gaps or mismatches. Repeated cycles of denaturation, hybridization and ligation amplify a short segment of DNA. LCR has also been used in combination with PCR to achieve enhanced detection of single-base changes. Segev, PCT Public. No. W09001069 A1 (1990). However, because the four oligonucleotides used in this assay can pair to form two short ligatable fragments, there is the potential for the generation of target-independent background signal. The use of LCR for mutant screening is limited to the examination of specific nucleic acid positions.\nSelf-Sustained Synthetic Reaction (3SR/NASBA)\nThe self-sustained sequence replication reaction (3SR) (Guatelli et al., Proc. Natl. Acad. Sci., 87:1874-1878 (1990), with an erratum at Proc. Natl. Acad. Sci., 87:7797 (1990)) is a transcription-based in vitro amplification system (Kwok et al., Proc. Natl. Acad. Sci., 86:1173-1177 (1989)) that can exponentially amplify RNA sequences at a uniform temperature. The amplified RNA can then be utilized for mutation detection (Fahy et al., PCR Meth. Appl., 1:25-33 (1991)). In this method, an oligonucleotide primer is used to add a phage RNA polymerase promoter to the 5′ end of the sequence of interest. In a cocktail of enzymes and substrates that includes a second primer, reverse transcriptase, RNase H, RNA polymerase and ribo- and deoxyribonucleoside triphosphates, the target sequence undergoes repeated rounds of transcription, cDNA synthesis and second-strand synthesis to amplify the area of interest. The use of 3SR to detect mutations is kinetically limited to screening small segments of DNA (e.g., 200-300 base pairs).\nQ-Beta (Qβ) Replicase\nIn this method, a probe which recognizes the sequence of interest is attached to the replicatable RNA template for Qβ replicase. A previously identified major problem with false positives resulting from the replication of unhybridized probes has been addressed through use of a sequence-specific ligation step. However, available thermostable DNA ligases are not effective on this RNA substrate, so the ligation must be performed by T4 DNA ligase at low temperatures (37° C.). This prevents the use of high temperature as a means of achieving specificity as in the LCR, the ligation event can be used to detect a mutation at the junction site, but not elsewhere.\nTable 1 below, lists some of the features desirable for systems useful in sensitive nucleic acid diagnostics, and summarizes the abilities of each of the major amplification methods (See also, Landgren, Trends in Genetics 9:199 (1993)).\nA successful diagnostic method must be very specific. A straight-forward method of controlling the specificity of nucleic acid hybridization is by controlling the temperature of the reaction. While the 3SR/NASBA, and Qβ systems are all able to generate a large quantity of signal, one or more of the enzymes involved in each cannot be used at high temperature (i.e., >55° C.). Therefore the reaction temperatures cannot be raised to prevent non-specific hybridization of the probes. If probes are shortened in order to make them melt more easily at low temperatures, the likelihood of having more than one perfect match in a complex genome increases. For these reasons, PCR and LCR currently dominate the research field in detection technologies.\nTABLE 1METHOD:PCR &3SRFEATUREPCRLCRLCRNASBAQβAmplifies Target++++Recognition of Independent+++++Sequences RequiredPerformed at High Temp.++Operates at Fixed Temp.++Exponential Amplification+++++Generic Signal Generation+Easily Automatable\nThe basis of the amplification procedure in the PCR and LCR is the fact that the products of one cycle become usable templates in all subsequent cycles, consequently doubling the population with each cycle. The final yield of any such doubling system can be expressed as: (1+X)n=y, where “X” is the mean efficiency (percent copied in each cycle), “n” is the number of cycles, and “y” is the overall efficiency, or yield of the reaction (Mullis, PCR Methods Applic., 1:1 (1991)). If every copy of a target DNA is utilized as a template in every cycle of a polymerase chain reaction, then the mean efficiency is 100%. If 20 cycles of PCR are performed, then the yield will be 220, or 1,048,576 copies of the starting material. If the reaction conditions reduce the mean efficiency to 85%, then the yield in those 20 cycles will be only 1.8520, or 220,513 copies of the starting material. In other words, a PCR running at 85% efficiency will yield only 21% as much final product, compared to a reaction running at 100% efficiency. A reaction that is reduced to 50% mean efficiency will yield less than 1% of the possible product.\nIn practice, routine polymerase chain reactions rarely achieve the theoretical maximum yield, and PCRs are usually run for more than 20 cycles to compensate for the lower yield. At 50% mean efficiency, it would take 34 cycles to achieve the million-fold amplification theoretically possible in 20, and at lower efficiencies, the number of cycles required becomes prohibitive. In addition, any background products that amplify with a better mean efficiency than the intended target will become the dominant products.\nAlso, many variables can influence the mean efficiency of PCR, including target DNA length and secondary structure, primer length and design, primer and dNTP concentrations, and buffer composition, to name but a few. Contamination of the reaction with exogenous DNA (e.g., DNA spilled onto lab surfaces) or cross-contamination is also a major consideration. Reaction conditions must be carefully optimized for each different primer pair and target sequence, and the process can take days, even for an experienced investigator. The laboriousness of this process, including numerous technical considerations and other factors, presents a significant drawback to using PCR in the clinical setting. Indeed, PCR has yet to penetrate the clinical market in a significant way. The same concerns arise with LCR, as LCR must also be optimized to use different oligonucleotide sequences for each target sequence. In addition, both methods require expensive equipment, capable of precise temperature cycling.\nMany applications of nucleic acid detection technologies, such as in studies of allelic variation, involve not only detection of a specific sequence in a complex background, but also the discrimination between sequences with few, or single, nucleotide differences. One method for the detection of allele-specific variants by PCR is based upon the fact that it is difficult for Taq polymerase to synthesize a DNA strand when there is a mismatch between the template strand and the 3′ end of the primer. An allele-specific variant may be detected by the use of a primer that is perfectly matched with only one of the possible alleles; the mismatch to the other allele acts to prevent the extension of the primer, thereby preventing the amplification of that sequence. This method has a substantial limitation in that the base composition of the mismatch influences the ability to prevent extension across the mismatch, and certain mismatches do not prevent extension or have only a minimal effect (Kwok et al., Nucl. Acids Res., 18:999 (1990)).\nA similar 3′-mismatch strategy is used with greater effect to prevent ligation in the LCR (Barany, PCR Meth. Applic., 1:5 (1991)). Any mismatch effectively blocks the action of the thermostable ligase, but LCR still has the drawback of target-independent background ligation products initiating the amplification. Moreover, the combination of PCR with subsequent LCR to identify the nucleotides at individual positions is also a clearly cumbersome proposition for the clinical laboratory.\nII. Direct Detection Technology\nWhen a sufficient amount of a nucleic acid to be detected is available, there are advantages to detecting that sequence directly, instead of making more copies of that target, (e.g., as in PCR and LCR). Most notably, a method that does not amplify the signal exponentially is more amenable to quantitative analysis. Even if the signal is enhanced by attaching multiple dyes to a single oligonucleotide, the correlation between the final signal intensity and amount of target is direct. Such a system has an additional advantage that the products of the reaction will not themselves promote further reaction, so contamination of lab surfaces by the products is not as much of a concern. Traditional methods of direct detection including Northern and Southern blotting and RNase protection assays usually require the use of radioactivity and are not amenable to automation. Recently devised techniques have sought to eliminate the use of radioactivity and/or improve the sensitivity in automatable formats. Two examples are the “Cycling Probe Reaction” (CPR), and “Branched DNA” (bDNA)\nThe cycling probe reaction (CPR) (Duck et al., BioTech., 9:142 (1990)), uses a long chimeric oligonucleotide in which a central portion is made of RNA while the two termini are made of DNA. Hybridization of the probe to a target DNA and exposure to a thermostable RNase H causes the RNA portion to be digested. This destabilizes the remaining DNA portions of the duplex, releasing the remainder of the probe from the target DNA and allowing another probe molecule to repeat the process. The signal, in the form of cleaved probe molecules, accumulates at a linear rate. While the repeating process increases the signal, the RNA portion of the oligonucleotide is vulnerable to RNases that may carried through sample preparation.\nBranched DNA (bDNA), described by Urdea et al., Gene 61:253-264 (1987), involves oligonucleotides with branched structures that allow each individual oligonucleotide to carry 35 to 40 labels (e.g., alkaline phosphatase enzymes). While this enhances the signal from a hybridization event, signal from non-specific binding is similarly increased.\nIII. Detection of Unknown Sequence Changes\nThe demand for tests which allow the detection of specific nucleic acid sequences and sequence changes is growing rapidly in clinical diagnostics. As nucleic acid sequence data for genes from humans and pathogenic organisms accumulates, the demand for fast, cost-effective, and easy-to-use tests for as yet unknown mutations within specific sequences is rapidly increasing.\nA handful of methods have been devised to scan nucleic acid segments for mutations. One option is to determine the entire gene sequence of each test sample (e.g., a bacterial isolate). For sequences under approximately 600 nucleotides, this may be accomplished using amplified material (e.g., PCR reaction products). This avoids the time and expense associated with cloning the segment of interest. However, specialized equipment and highly trained personnel are required, and the method is too labor-intense and expensive to be practical and effective in the clinical setting.\nIn view of the difficulties associated with sequencing, a given segment of nucleic acid may be characterized on several other levels. At the lowest resolution, the size of the molecule can be determined by electrophoresis by comparison to a known standard run on the same gel. A more detailed picture of the molecule may be achieved by cleavage with combinations of restriction enzymes prior to electrophoresis, to allow construction of an ordered map. The presence of specific sequences within the fragment can be detected by hybridization of a labeled probe, or the precise nucleotide sequence can be determined by partial chemical degradation or by primer extension in the presence of chain-terminating nucleotide analogs.\nFor detection of single-base differences between like sequences, the requirements of the analysis are often at the highest level of resolution. For cases in which the position of the nucleotide in question is known in advance, several methods have been developed for examining single base changes without direct sequencing. For example, if a mutation of interest happens to fall within a restriction recognition sequence, a change in the pattern of digestion can be used as a diagnostic tool (e.g., restriction fragment length polymorphism (RFLP) analysis).\nSingle point mutations have been also detected by the creation or destruction of RFLPs. Mutations are detected and localized by the presence and size of the RNA fragments generated by cleavage at the mismatches. Single nucleotide mismatches in DNA heteroduplexes are also recognized and cleaved by some chemicals, providing an alternative strategy to detect single base substitutions, generically named the “Mismatch Chemical Cleavage” (MCC) (Gogos et al., Nucl. Acids Res., 18:6807-6817 (1990)). However, this method requires the use of osmium tetroxide and piperidine, two highly noxious chemicals which are not suited for use in a clinical laboratory.\nRFLP analysis suffers from low sensitivity and requires a large amount of sample. When RFLP analysis is used for the detection of point mutations, it is, by its nature, limited to the detection of only those single base changes which fall within a restriction sequence of a known restriction endonuclease. Moreover, the majority of the available enzymes have 4 to 6 base-pair recognition sequences, and cleave too frequently for many large-scale DNA manipulations (Eckstein and Lilley (eds.), Nucleic Acids and Molecular Biology, vol. 2, Springer-Verlag, Heidelberg (1988)). Thus, it is applicable only in a small fraction of cases, as most mutations do not fall within such sites.\nA handful of rare-cutting restriction enzymes with 8 base-pair specificities have been isolated and these are widely used in genetic mapping, but these enzymes are few in number, are limited to the recognition of G+C-rich sequences, and cleave at sites that tend to be highly clustered (Barlow and Lehrach, Trends Genet., 3:167 (1987)). Recently, endonucleases encoded by group I introns have been discovered that might have greater than 12 base-pair specificity (Perlman and Butow, Science 246:1106 (1989)), but again, these are few in number.\nIf the change is not in a recognition sequence, then allele-specific oligonucleotides (ASOs), can be designed to hybridize in proximity to the unknown nucleotide, such that a primer extension or ligation event can be used as the indicator of a match or a mis-match. Hybridization with radioactively labeled allelic specific oligonucleotides (ASO) also has been applied to the detection of specific point mutations (Conner et al., Proc. Natl. Acad. Sci., 80:278-282 (1983)). The method is based on the differences in the melting temperature of short DNA fragments differing by a single nucleotide. Stringent hybridization and washing conditions can differentiate between mutant and wild-type alleles. The ASO approach applied to PCR products also has been extensively utilized by various researchers to detect and characterize point mutations in ras genes (Vogelstein et al., N. Eng. J. Med., 319:525-532 (1988); and Farr et al., Proc. Natl. Acad. Sci., 85:1629-1633 (1988)), and gsp/gip oncogenes (Lyons et al., Science 249:655-659 (1990)). Because of the presence of various nucleotide changes in multiple positions, the ASO method requires the use of many oligonucleotides to cover all possible oncogenic mutations.\nWith either of the techniques described above (i.e., RFLP and ASO), the precise location of the suspected mutation must be known in advance of the test. That is to say, they are inapplicable when one needs to detect the presence of a mutation of an unknown character and position within a gene or sequence of interest.\nTwo other methods rely on detecting changes in electrophoretic mobility in response to minor sequence changes. One of these methods, termed “Denaturing Gradient Gel Electrophoresis” (DGGE) is based on the observation that slightly different sequences will display different patterns of local melting when electrophoretically resolved on a gradient gel. In this manner, variants can be distinguished, as differences in melting properties of homoduplexes versus heteroduplexes differing in a single nucleotide can detect the presence of mutations in the target sequences because of the corresponding changes in their electrophoretic mobilities. The fragments to be analyzed, usually PCR products, are “clamped” at one end by a long stretch of G-C base pairs (30-80) to allow complete denaturation of the sequence of interest without complete dissociation of the strands. The attachment of a GC “clamp” to the DNA fragments increases the fraction of mutations that can be recognized by DGGE (Abrams et al., Genomics 7:463-475 (1990)). Attaching a GC clamp to one primer is critical to ensure that the amplified sequence has a low dissociation temperature (Sheffield et al., Proc. Natl. Acad. Sci., 86:232-236 (1989); and Lerman and Silverstein, Meth. Enzymol., 155:482-501 (1987)). Modifications of the technique have been developed, using temperature gradients (Wartell et al., Nucl. Acids Res., 18:2699-2701 (1990)), and the method can be also applied to RNA:RNA duplexes (Smith et al., Genomics 3:217-223 (1988)).\nLimitations on the utility of DGGE include the requirement that the denaturing conditions must be optimized for each type of DNA to be tested. Furthermore, the method requires specialized equipment to prepare the gels and maintain the needed high temperatures during electrophoresis. The expense associated with the synthesis of the clamping tail on one oligonucleotide for each sequence to be tested is also a major consideration. In addition, long running times are required for DGGE. The long running time of DGGE was shortened in a modification of DGGE called constant denaturant gel electrophoresis (CDGE) (Borrensen et al., Proc. Natl. Acad. Sci. USA 88:8405 (1991)). CDGE requires that gels be performed under different denaturant conditions in order to reach high efficiency for the detection of unknown mutations.\nAn technique analogous to DGGE, termed temperature gradient gel electrophoresis (TGGE), uses a thermal gradient rather than a chemcial denaturant gradient (Scholz, et al., Hum. Mol. Genet. 2:2155 (1993)). TGGE requires the use of specialized equipment which can generate a temperature gradient perpendicularly oriented relative to the electrical field. TGGE can detect mutations in relatively small fragments of DNA therefore scanning of large gene segments requires the use of multiple PCR products prior to running the gel.\nAnother common method, called “Single-Strand Conformation Polymorphism” (SSCP) was developed by Hayashi, Sekya and colleagues (reviewed by Hayashi, PCR Meth. Appl., 1:34-38, (1991)) and is based on the observation that single strands of nucleic acid can take on characteristic conformations in non-denaturing conditions, and these conformations influence electrophoretic mobility. The complementary strands assume sufficiently different structures that one strand may be resolved from the other. Changes in sequences within the fragment will also change the conformation, consequently altering the mobility and allowing this to be used as an assay for sequence variations (Orita, et al., Genomics 5:874-879, (1989)).\nThe SSCP process involves denaturing a DNA segment (e.g., a PCR product) that is labelled on both strands, followed by slow electrophoretic separation on a non-denaturing polyacrylamide gel, so that intra-molecular interactions can form and not be disturbed during the run. This technique is extremely sensitive to variations in gel composition and temperature. A serious limitation of this method is the relative difficulty encountered in comparing data generated in different laboratories, under apparently similar conditions.\nThe dideoxy fingerprinting (ddF) is another technique developed to scan genes for the presence of unknown mutations (Liu and Sommer, PCR Methods Appli., 4:97 (1994)). The ddF technique combines components of Sanger dideoxy sequencing with SSCP. A dideoxy sequencing reaction is performed using one dideoxy terminator and then the reaction products are electrophoresised on nondenaturing polyacrylamide gels to detect alterations in mobility of the termination segments as in SSCP analysis. While ddF is an improvement over SSCP in terms of increased sensitivity, ddF requires the use of expensive dideoxynucleotides and this technique is still limited to the analysis of fragments of the size suitable for SSCP (i.e., fragments of 200-300 bases for optimal detection of mutations).\nIn addition to the above limitations, all of these methods are limited as to the size of the nucleic acid fragment that can be analyzed. For the direct sequencing approach, sequences of greater than 600 base pairs require cloning, with the consequent delays and expense of either deletion sub-cloning or primer walking, in order to cover the entire fragment. SSCP and DGGE have even more severe size limitations. Because of reduced sensitivity to sequence changes, these methods are not considered suitable for larger fragments. Although SSCP is reportedly able to detect 90% of single-base substitutions within a 200 base-pair fragment, the detection drops to less than 50% for 400 base pair fragments. Similarly, the sensitivity of DGGE decreases as the length of the fragment reaches 500 base-pairs. The ddF technique, as a combination of direct sequencing and SSCP, is also limited by the relatively small size of the DNA that can be screened.\nClearly, there remains a need for a method that is less sensitive to size so that entire genes, rather than gene fragments, may be analyzed. Such a tool must also be robust, so that data from different labs, generated by researchers of diverse backgrounds and skills will be comparable. Ideally, such a method would be compatible with “multiplexing,” (i.e., the simultaneous analysis of several molecules or genes in a single reaction or gel lane, usually resolved from each other by differential labelling or probing). Such an analytical procedure would facilitate the use of internal standards for subsequent analysis and data comparison, and increase the productivity of personnel and equipment. The ideal method would also be easily automatable."}
{"text": "Mobile applications typically have a compulsory/optional step for users to register that involves entering email addresses and mobile numbers. For some applications, registration is an obvious requirement for accessing certain functionality. For example, banking applications may be able to locate the nearest ATM without registration, but to display account information, users must be registered. While application developers wish to keep the registration process simple, some applications may make customers uncomfortable with a process that is too short because the process does not seem secure. As such, application developers attempt to provide an easy registration process while also ensuring users that the registration process is secure.\nOver time, mobile device users may download and register with many applications. Typically, a user's many applications are tied to the user's mobile device/phone number. The user, however, may be required to change mobile numbers due to geographical location changes or for personal reasons. When a user changes his/her number, the user usually prefers to update the new number with the individual services or initiate a request to respective authorities. Unfortunately, the user may find it difficult to recall which services he/she needs to update before deactivating the user's old mobile number."}
{"text": "1. Field of the Invention\nThe invention relates to a method of manufacturing a plurality of electronic multilayer components, each of which comprises alternately stacked electrically conductive and electrically insulating layers, the electrically conductive layers being electrically connected in a periodically alternate arrangement to different edges of the multilayer component. Such components may receive application as multilayer capacitors or multilayer actuators, example.\n2. Discussion of the Related Art\nA method as described in the opening paragraph is known from U.S. Pat. No. 3,326,718, in which layers of electrically conductive and electrically insulating material are alternately deposited onto a flat substrate through an apertured mask, the planes of the substrate and mask being mutually parallel. In the case of the insulating material, the depository flux is directed at right angles to the plane of the mask, so that it passes through the aperture in a perpendicular direction. However, in the case of the conductive material, the depository flux is directed through the aperture at a non-perpendicular angle .alpha. with respect to the substrate surface. Moreover, although consecutive conductive layers are deposited using the same value of .alpha., the depository fluxes for such consecutive layers are not mutually parallel, but instead arise from sources located at diametrically opposite sides of the aperture. As a result, consecutive conductive layers demonstrate only a partial mutual overlap, as illustrated in FIG. 2 of the cited U.S. patent. At the same time, as shown in FIG. 4 of that patent, conductive layers (56, 56') having an odd ordinal number make mutual electrical contact at one side (62) of the component, and conductive layers (68, 68') having an even ordinal number make mutual electrical contact at an other side (76) of the component.\nThe known method has a number of disadvantages. In particular, the number of conductive layers which can be deposited in this manner is severely limited. This is because, as the stack of layers on the substrate increases in height, that stack will itself begin to partially eclipse the depository fluxes of conductive material, and will eventually prevent the desired mutual contact of every second conductive layer at the edge of the component. In such a scenario, the finished component will have to be provided along its sides with blanketing layers of conductive material (such as solder), for the purpose of achieving uninterrupted interconnection of the conductive layers terminating at each given side. However, because such blanketing layers are at the sides of the component, they are not directly compatible with surface mounting techniques, which require the component's electrical contacts to be located in a single plane."}
{"text": "The present invention relates to compounds, processes for their preparation, compositions containing them, and their therapeutic use as vaccine adjuvants and in the treatment of various disorders.\nThe innate immune system recognises microbes via a limited number of germline-encoded Pattern-Recognition Receptors (PRRs) which have a number of important characteristics.\nToll-like receptors (TLR) are a family of structurally related PRRs that detect highly conserved microbial components common to large classes of pathogens. TLRs are expressed on immune cells and upon activation mobilize defense mechanisms aimed at eliminating the invading pathogens. Of the more than ten known TLRs that have been identified in humans, some appear to be restricted to cytoplasmic compartments and involved in the detection of non-self nucleic acids (TLRs 3, 7, 8, and 9). See, e.g., Akira et al., Nat Rev Immunol 2004, 4, 499-511; O'Neill, et al., Nat Rev Immunol 2013, 13, 453-460.\nActivation of TLRs regulates intracellular signaling pathways leading to the expression of inflammatory cytokines/chemokines and type I interferons (IFNα/β), which can lead to the preferential enhancement of antigen-specific humoral and cell-mediated immune responses.\nTLR7 and TLR8 are members of the subgroup of TLRs (TLRs 3, 7, 8, and 9) localised in the endosomal compartment of cells. TLR7 plays a key rôle in anti-viral defence via the recognition of ssRNA (Diebold S.S. et al, Science, 2004: 303, 1529-1531; and Lund J. M. et al, PNAS, 2004: 101, 5598-5603). TLR7 has a restricted expression-profile in human and is expressed predominantly by B cells and plasmacytoid dendritic cells (pDC), and to a lesser extent by monocytes. Plasmacytoid DCs are a unique population of lymphoid-derived dendritic cells (typically 0.2-0.8% of Peripheral Blood Mononuclear Cells (PBMCs)) and are the primary type I interferon-producing cells secreting high levels of interferon-alpha (IFNα) and interferon-beta (IFNβ) in response to viral infections (Liu Y-J, Annu. Rev. Immunol., 2005: 23, 275-306).\nSmall-molecule agonists of TLR7 have been described which can induce cytokines in animals and in man (Takeda K. et al, Annu. Rev. Immunol., 2003: 21, 335-76). TLR7 agonists include imidazoquinoline compounds such as imiquimod and resiquimod, oxoadenine analogues and also nucleoside analogues such as loxoribine and 7-thia-8-oxoguanosine, which are known to induce interferon alpha. International Patent Application publication number WO 2007/034882 (PCT/JP2006/318758; Dainippon Sumitomo Pharma Co. Ltd./AstraZeneca Aktiebolag) discloses certain adenine compounds identified as useful as medicine.\nCertain adenine derivative compounds have been shown to be inducers of human interferon. Compounds which induce human interferon may be useful as vaccine adjuvants, as well as in the treatment of various disorders, including infectious diseases, asthma, cancer, inflammatory conditions, and allergic diseases. It is thus desirable to provide compounds having selectivity and/or potency for TLR7/8 and high relative cytokine induction."}
{"text": "1. Field of the Invention\nThe invention relates to devices for aligning articles, and in particular, devices for aligning sheets or the like into edge-aligned stacks.\n2. Prior Art\nIn handling planar articles such as paper sheets outputted from printers, presses and electrophotographic copiers, it is often required to stack the sheets into aligned stacks for operations such as cutting, stapling and binding. The process of forming stacks of aligned sheets may be done by mechanical means or nonmechanical means.\nWhen sheets are properly aligned with one another by means of the human hands, the process is often referred to as nonmechanical. The partially-aligned stack of sheets is held by the hand. By tapping alternately adjacent edges of the sheet against a flat surface, forces the sheets into alignment. Although this procedure works satisfactorily, it is not well suited for commercial adaptation. Moreover, the procedure is time-consuming and expensive.\nIn an attempt to circumvent some of the disadvantages associated with human hands, mechanical devices have been used to align sheets. One type of prior art mechanical aligner consists of an inclined table with a pair of jogger arms pivotally mounted to the table. The jogger arms have a pair of paddle portions extending upwardly above the level of the table along two adjacent sides. The table is inclined towards the paddle portions of the jogger arm. As such, a gravitational force is imparted to the sheets along the direction of the incline. The force helps to bring the sheets into alignment. A driving means consisting of a motor-driven camming system activates the jogger arms which causes pivotal movement of the paddles. The paddles tap against the sides or edges of the sheets delivered on the table to form a properly aligned stack. A more detailed description of the above prior art joggers is given in U.S. Pat. No. 3,593,992.\nAnother type of prior art aligners is described in U.S. Pat. No. 3,083,014. In that patent, sheet-like articles to be formed into edge-aligned stacks are delivered serially to a stacker and jogger mechanism. The trailing end and leading end of the sheets are overlapped. The stacker and jogger mechanism consists of an alignment surface and a movable table for supporting the sheet articles. Two pair of resilient bladed rotating paddle wheels are mounted relative to the table. One pair of the paddle wheels is disposed on opposite edges or sides of the table. The paddle wheels in each pair are in spaced relation on its respective side of the table. The paddles are inclined with respect to the table. As sheet-like articles are delivered to the table, in the direction of paddle rotation, the rotating resilient paddle wheels contact and lightly impact the opposite edges of the sheets to impart a jogging or vibratory action which aligns the sheet-like articles against the alignment surface.\nAlthough the above-described aligners probably work well for their intended purpose, there are times when the above aligners do not align the sheets with sufficient accuracy. For example, if some of the sheets in a particular size classification (such as 81/2\".times.14\" etc.) are slightly less than the stated size for that classification. The inability of the aligners to accurately align sheets in a stack wherein the dimension of some sheets are slightly less than the stated dimension stems from the fact that the prior art aligners all work on the edges of the sheets. The smaller sheets in a mixed size stack do not extend to the edges of the stack, therefore, tapping on the side of the stack does not align the sheets since there is no contact between the tapping element and the smaller size sheets."}
{"text": "The present invention relates to a ski boot closure device.\nSeveral devices for ski boots are currently known which allow for example to close the quarters and lock one or more pressers arranged inside the shell.\nA U.S. Pat. No. 4,571,855, is known, for example, to disclose a lever for the closure of the quarters which uses two cables and has a system for the differentiated tensioning thereof.\nThis differentiation is intrinsically necessary, since when the quarters are opened the working length of the cables must be different, and this has made necessary to provide a guide member for said cables above said lever at the rear quarter.\nThis undoubtedly produces a disadvantage, because too much space is taken up on the rear quarter and, besides the undesirable aesthetic aspect, no other devices can be placed at the rear quarter.\nA French patent, No. 84.16966 filed on Oct. 30, 1984, is also known which describes a rear-entry ski boot comprising a lever provided with two coaxial winders for two separate cables.\nA first cable closes the quarters and the other cable operates a foot instep presser.\nSaid two separate cables can be actuated by means of a single knob, and means are provided in said knob adapted to provide different transmission ratios; the different transmission ratios are required both to obtain the different take-up of the cables and to balance the efforts for the tightening thereof.\nSaid known device, however, is relatively complicated structurally, because it requires a different degree of take-up of the cables for an equal number of turns imparted to the single knob.\nClosure and adjustment devices are also known which are constituted by traction elements associated with a vertical lever so that an increase in the tension of said traction elements is followed by a shift of the point of application of the resisting force, and thus the effort on the lever is variable.\nOther known devices comprise a vertical lever associated for example with the rear quarter of the boot and sliders axially movable on a threaded stem associated with a lever for the take-up of the cables. When the lever is closed, the slider and stem are subjected to a moment that, after some time, eventually leads to the stripping of the thread of the threaded stem or of the complementary thread of the sliders."}
{"text": "Membrane-derived microvesicles (MVs) include a range of extracellular vesicles, including exosomes, microparticles and shed MVs secreted by many cell types under both normal physiological and pathological conditions [1]. As intercellular communication tools, MVs have been reported to have roles in a wide range of cellular functions: immunological modulation, coagulation, and tumor progression, including angiogenesis and metastasis [2,3]. Additionally, they have been reported to serve as vehicles for transferring cargo (mRNA, miRNA, non-coding RNAs, proteins and oncogenes) between cells [4-8]. The mRNA content in MVs opens new research opportunities from cancer diagnostics to gene therapy applications [9-12]. Despite the intensive research in analyzing the RNA content of MVs, it remains unclear how RNAs are directed to MVs.\nMechanistically, cis-acting regulatory sequences and trans-acting proteins are considered as the main driving forces of mRNA localization within cells and have an important role in post-transcriptional regulation. Such sequences, also known as zip codes, are typically found in the 3′-untranslated regions (3′ UTRs) of mRNA transcripts and mediate binding of a ribonuclear protein complex to the mRNA which temporarily blocks its translation and mediates movement via the cytoskeleton to a cellular location where it releases the mRNA and translation commences [13,14].\nIn addition to zip codes, microRNAs (miRNAs) also have a critical role in post-transcriptional regulation of mRNAs. miRNAs are small, non-coding, single-stranded RNA molecules [˜21-23 nucleotides (nt) long] that regulate levels of gene expression in many organisms [15]. miRNAs mediate post-transcriptional regulation in three ways, by mRNA degradation, mRNA destabilization via deadenylation and translational repression [16]. In addition to negative regulation, miRNAs have also been reported to function in the activation of translation in some cases [17]. Although zip codes and miRNA target sequences are both found in the 3′UTRs of the mRNA transcripts, there has been no report of cooperation between these two regulatory mechanisms in mRNA fate."}
{"text": "The production of light olefins (ethylene, propylene and butenes) from various hydrocarbon feedstocks utilizes the technique of pyrolysis, or steam cracking. Pyrolysis involves heating the feedstock sufficiently to cause thermal decomposition of the larger molecules. The pyrolysis process, however, produces molecules which tend to combine to form high molecular weight materials known as tars. Tars are high-boiling point, viscous, reactive materials that can foul equipment under certain conditions.\nThe formation of tars, after the pyrolysis effluent leaves the steam cracking furnace can be minimized by rapidly reducing the temperature of the effluent exiting the pyrolysis unit to a level at which the tar-forming reactions are greatly slowed.\nOne technique used to cool pyrolysis unit effluent and remove the resulting heavy oils and tars employs heat exchangers followed by a water quench tower in which the condensibles are removed. This technique has proven effective when cracking light gases, primarily ethane, propane and butane, because crackers that process light feeds, collectively referred to as gas crackers, produce relatively small quantities of tar. As a result, heat exchangers can efficiently recover most of the valuable heat without fouling and the relatively small amount of tar can be separated from the water quench albeit with some difficulty.\nThis technique is, however, not satisfactory for use with steam crackers that crack naphthas and heavier feedstocks, collectively referred to as liquid crackers, since liquid crackers generate much larger quantities of tar than gas crackers. Heat exchangers can be used to remove some of the heat from liquid cracking, but only down to the temperature at which tar begins to condense. Below this temperature, conventional heat exchangers cannot be used because they would foul rapidly from accumulation and thermal degradation of tar on the heat exchanger surfaces. In addition, when the pyrolysis effluent from these feedstocks is quenched, some of the heavy oils and tars produced have approximately the same density as water and can form stable oil/water emulsions. Moreover, the larger quantity of heavy oils and tars produced by liquid cracking would render water quench operations ineffective, making it difficult to raise steam from the condensed water and to dispose of excess quench water and the heavy oil and tar in an environmentally acceptable manner.\nAccordingly in most commercial liquid crackers, cooling of the effluent from the cracking furnace is normally achieved using a system of transfer line heat exchangers, a primary fractionator, and a water quench tower or indirect condenser. For a typical naphtha feedstock, the transfer line heat exchangers cool the process stream to about 700° F. (370° C.), efficiently generating super-high pressure steam that can be used elsewhere in the process. The primary fractionator is normally used to condense and separate the tar from the lighter liquid fraction, known as pyrolysis gasoline, and to recover the heat between about 700° F. (370° C.) and about 200° F. (90° C.). The water quench tower or indirect condenser further cools the gas stream exiting the primary fractionator to about 104° F. (40° C.) to condense the bulk of the dilution steam present and to separate pyrolysis gasoline from the gaseous olefinic product, which is then sent to a compressor.\nThe primary fractionator, however, is a very complex piece of equipment that typically includes an oil quench section, a primary fractionator tower and one or more external oil pumparound loops. At the quench section, quench oil is added to cool the effluent stream to about 400 to 6504° F. (200 to 343° C.), thereby condensing tar present in the stream. In the primary fractionator tower, the condensed tar is separated from the remainder of the stream, heat is removed in one or more pumparound zones by circulating oil and a pyrolysis gasoline fraction is separated from heavier material in one or more distillation zones. In the one or more external pumparound loops, oil, which is withdrawn from the primary fractionator, is cooled using indirect heat exchangers and then returned to the primary fractionator or the direct quench point.\nThe primary fractionator with its associated pumparounds is the most expensive component in the entire cracking system. The primary fractionator tower itself is the largest single piece of equipment in the process, typically being about twenty-five feet in diameter and over a hundred feet high for a medium size liquid cracker. The tower is large because it is in effect fractionating two minor components, tar and pyrolysis gasoline, in the presence of a large volume of low pressure gas. The pumparound loops are likewise large, handling over 3 million pounds per hour of circulating oil in the case of a medium size cracker. Heat exchangers in the pumparound circuit are necessarily large because of high flow rates, close temperature approaches needed to recover the heat at useful levels, and allowances for fouling.\nIn addition, the primary fractionator has a number of other limitations and problems. In particular, heat transfer takes place twice, i.e., from the gas to the pumparound liquid inside the tower and then from the pumparound liquid to the external cooling service. This effectively requires investment in two heat exchange systems, and imposes two temperature approaches (or differentials) on the removal of heat, thereby reducing thermal efficiency.\nMoreover, despite the fractionation that takes place between the tar and gasoline streams, both streams often need to be processed further. Sometimes the tar needs to be stripped to remove light components, whereas the gasoline may need to be refractionated to meet its end point specification.\nFurther, the primary fractionator tower and its pumparounds are prone to fouling. Coke accumulates in the bottom section of the tower and must eventually be removed during plant turnarounds. The pumparound loops are also subject to fouling, requiring removal of coke from filters and periodic cleaning of fouled heat exchangers. Trays and packing in the tower are sometimes subject to fouling, potentially limiting plant production. The system also contains a significant inventory of flammable liquid hydrocarbons, which is not desirable from an inherent safety standpoint.\nThere is therefore a need for a simplified method for cooling pyrolysis unit effluent and removing the resulting heavy oils and tars which obviates the need for a primary fractionator tower and its ancillary equipment.\nU.S. Pat. Nos. 4,279,733 and 4,279,734 propose cracking methods using a quencher, indirect heat exchanger and fractionator to cool effluent, resulting from steam cracking.\nU.S. Pat. Nos. 4,150,716 and 4,233,137 propose a heat recovery apparatus comprising a pre-cooling zone where the effluent resulting from steam cracking is brought into contact with a sprayed quenching oil, a heat recovery zone, and a separating zone.\nLohr et al., “Steam-cracker Economy Keyed to Quenching,” Oil & Gas Journal, Vol. 76 (No. 20), pp. 63-68, (1978), proposes a two-stage quenching involving indirect quenching with a transfer line heat exchanger to produce high-pressure steam along with direct quenching with a quench oil to produce medium-pressure steam.\nU.S. Pat. Nos. 5,092,981 and 5,324,486 propose a two-stage quench process for effluent resulting from steam cracking furnace comprising a primary transfer line exchanger which functions to rapidly cool furnace effluent and to generate high temperature steam and a secondary transfer line exchanger which functions to cool the furnace effluent to as low a temperature as possible consistent with efficient primary fractionator or quench tower performance and to generate medium to low pressure steam.\nU.S. Pat. No. 5,107,921 proposes transfer line exchangers having multiple tube passes of different tube diameters. U.S. Pat. No. 4,457,364 proposes a close-coupled transfer line heat exchanger unit.\nU.S. Pat. No. 3,923,921 proposes a naphtha steam cracking process comprising passing effluent through a transfer line exchanger to cool the effluent and thereafter through a quench tower.\nWO 93/12200 proposes a method for quenching the gaseous effluent from a hydrocarbon pyrolysis unit by passing the effluent through transfer line exchangers and then quenching the effluent with liquid water so that the effluent is cooled to a temperature in the range of 220° F. to 266° F. (105° C. to 130° C.), such that heavy oils and tars condense, as the effluent enters a primary separation vessel. The condensed oils and tars are separated from the gaseous effluent in the primary separation vessel and the remaining gaseous effluent is passed to a quench tower where the temperature of the effluent is reduced to a level at which the effluent is chemically stable.\nEP 205 205 proposes a method for cooling a fluid such as a cracked reaction product by using transfer line exchangers having two or more separate heat exchanging sections.\nU.S. Pat. No. 5,294,347 proposes that in ethylene manufacturing plants, a water quench column cools gas leaving a primary fractionator and that in many plants, a primary fractionator is not used and the feed to the water quench column is directly from a transfer line exchanger.\nJP 2001-40366 proposes cooling mixed gas in a high temperature range with a horizontal heat exchanger and then with a vertical heat exchanger having its heat exchange planes installed in the vertical direction. A heavy component condensed in the vertical exchanger is thereafter separated by distillation at downstream refining steps.\nWO 00/56841; GB 1,390,382; GB 1,309,309; and U.S. Pat. Nos. 4,444,697; 4,446,003; 4,121,908; 4,150,716; 4,233,137; 3,923,921; 3,907,661; and 3,959,420; propose various apparatus for quenching a hot cracked gaseous stream wherein the hot gaseous stream is passed through a quench pipe or quench tube wherein a liquid coolant (quench oil) is injected."}
{"text": "Field of the Invention\nThe present invention relates to an image generation apparatus, a control method therefor, and a computer-readable storage medium, and especially relates to a technique to generate a privacy protection image.\nDescription of the Related Art\nAlong with the widespread installation of surveillance cameras, protection of privacy of individuals shown on the video of the surveillance cameras is becoming increasingly important Japanese Patent Laid-Open No. 2009-225398 describes a technique to protect privacy by generating a foreground image and a background image, which respectively represent an image of people and show no people, from an input image, and making an image, resulting from applying mask processing and filter processing to the foreground image, overlap the background image.\nJapanese Patent Laid-Open No. 2009-225398 is based on the premise that a background image is generated in a state where there is no person targeted for protection. If a background image is generated from the video of a surveillance camera in a state where a person exists, the background image will include an image of that person. That is to say, the image of that person is distinguished as a background, and is not generated as a foreground image. As a result, the privacy of that person may not be protected. Furthermore, if a person stays motionless for a long time from, for example, sitting on a chair, a person area may be reflected in a background image as time passes. In this case also, that person area is not treated as an area targeted for protection, and the privacy of that person cannot be protected."}
{"text": "The present invention relates generally to welding machines. More specifically, it relates to a welding machine that senses one or more conditions of a welding job and automatically sets operating parameters.\nMetal-inert-gas (MIG) welding is one well known type of welding process. There are a number of different types of power supplies used for MIG welding. Examples of MIG power supplies include phase controlled, pulse width modulated and inverter power supplies. MIG welding is performed using a variety of wire diameters, wire material, and gasses. The wire and gas used depends on the plate or workpiece thickness and material. Wire diameter, wire material, gas and plate thickness are referred to herein as welding conditions.\nMost MIG welding machines use at least three parameters to control the welding process. These parameters are: wire feed speed and/or current, voltage, and inductance. Some MIG welding machines also have an inductance control that affects the response of the power source or supply. Typically, a MIG welding machine will have switches and/or knobs on the front panel to select some or all of the operating parameters. The process of setting one or more of these parameters is called herein machine set up.\nSome prior art welding machines (as used herein welding machine refers to the power source, controller, wire feeder, and may include a gas source and other ancillary equipment used to effectuate a weld) are capable of controlling the welding arc to optimize the weld. However, proper control requires optimal operating parameters (current and/or wire feed speed and voltage) for the particular welding conditions (wire diameter, wire material and gas mixture used, as well as the plate thickness and joint type that is being welded). Prior art welding machines require the operator to calculate setup parameters (i.e., voltage, current and/or wire feed speed) from tables or equations using the above welding conditions inputs as independent variables. Alternatively, the appropriate settings may be chosen based on past experience, or by trial and error.\nIf the operator provides erroneous data, or doesn\"\"t properly calculate the setup parameters, the result may be poor quality welds or inefficient use of the welding machine or consumables (gas and wire). Thus, the weld quality is dependent upon the operator properly determining setup parameters. It is not unusual for the operators to improperly select the setup parameters and have such poor quality welds that a service call is required.\nAccordingly, it would be desirable for a welding machine to automatically sense one or more of the conditions which enter into the decision for setting wire feed speed and/or current and voltage. Moreover, it would be desirable to automatically set one or more of feed speed and/or current or voltage, in response to the sensed condition. Such a welding machine should be capable of utilizing the wide variety of power supplies.\nAccording to a first embodiment of the invention a welding machine includes a power source and a wire feeder. At least one sensor senses at least one welding condition. A set up circuit is coupled to the sensor, the power source, and the wire feeder. The set up circuit automatically sets up the welding machine in response to the sensor.\nAccording to one alternative the sensor is a wire sensor. In another alternative the welding machine includes a source of gas and a gas sensor. Another alternative includes a workpiece thickness sensor. The set up circuit further automatically sets up the welding machine in response to the additional sensors.\nYet another embodiment uses a diameter sensor and a material sensor as part of the wire sensor. Other alternative configurations of the sensor include a bar code reader, an analog proximity detector, a plurality of proximity detector, and a displacement sensor.\nOther embodiments provides for the gas sensor to include a cyclic voltametry gas microsensor, a bar code reader, or a plurality of gas connections. The thickness detector may include a strain gauge or displacement sensor.\nOther principal features and advantages of the invention will become apparent to those skilled in the art upon review of the following drawings, the detailed description and the appended claims."}
{"text": "A. Field of the Invention\nThe present invention relates to a conduit comprising 2 or more pipes that are twisted together. Said conduit is an integral part of a geothermal heating and cooling system. Geothermal heating and cooling systems are known for their superior performance in delivering efficient heating and cooling to homes, industrial buildings and residential and industrial complexes, as well as being environmentally clean and cost effective. See, for example,    http://www.igshpa.okstate.edu/geothermal/geothermal.htmI;    www.summitmechsystems.com/pages/3.1.html;    www.renewableheating101.com/geothermal/loops;    http://minnesotageothermalheatpumpassociation.com/geothermal/earth-loop-options/; and http://www.informedbuilding.com/Geothermal/Main16/Types-of-Geotherm.However, a barrier to the wide spread use of geothermal heating and cooling systems is the high cost of installation of the ground loop of pipes that the system requires. Also, the presently available ground-loop pipes are not ideal in terms of heat transfer and utilization of the borehole required to install the vertical ground loop pipes.\nB. Description of the Related Art\nThe art has attempted to overcome these barriers to market entry for this efficient heating and cooling system. The aforementioned websites discuss the currently available ground loop technology. For example, the commonly used ground loop technologies are: horizontal ground loops, vertical ground loops, and slinky coil ground loops. However, horizontal ground loops require a substantial amount of land. Currently available vertical loops, including multiple pipe vertical loops use less land, but their configuration does not optimize heat transfer, as does the present invention. Finally, the slinky coil ground loop is a variation of the horizontal ground loop and it too requires a substantial amount of land.\nAlso, inventors have sought patents on conduits to solve the problems with the current technology. For example, U.S. Pat. No. 5,630,447 ('447) discloses a pipe design that utilizes the entire borehole; and therefore transfers more heat to the ground. Further, the '447 invention allows for reduction in the size of the borehole required for a pipe capable of handling a specified flow of heat transfer fluid. The '447 invention, however, has some limitations. Said limitation being the pipe design and the cost of said design both in terms of time and money. The standards set by ASTM specify that pressure pipes have a round configuration. The pressure rating is derived by a combination of material strength and diameter to pipe wall thickness ration. The smaller this ration the higher the pressure rating of the pipe will be. The pipe represented in U.S. Pat. No. 5,630,447 is not round according to the standard, set by ASTM, to determine pressure tolerance of a pipe, and therefore cannot be governed by the same standard. Thus, new standards will need to be written and approved by standard setting bodies such as ASTM and IGSHPA. This process could be costly and time consuming. Additionally the U.S. '447 pipe inherently keeps the heat transferring fluid in the in and out flow pipes in close proximity to each other; thus causing heat contamination from the inflow to the outflow pipe. U.S. Pat. No. 5,630,447 does address this problem by introducing the notion of an insulating space between the inflow and outflow pipe. However, this design further complicates the pipe design certification issues.\nFinally, U.S. Pat. No. 5,477,914 ('914) discloses a ground source heat exchanger unit comprising a primary conduit and a plurality of secondary conduits for receiving heat transfer fluid. Said secondary conduits are spaced apart from each other. The '914 system is not designed for optimal use of the borehole due to the spacing between the secondary conduits. FIG. 1 of the '914 disclosure illustrates the fact that the '914 system requires greater land usage than a typical narrow borehole installation. Since the borehole is a very costly part of the installation of these systems, the '914 design becomes costly to install because of the larger diameter borehole required by the '914 system.\nThe '914 system will be inherently more cumbersome to manage because of the flexibility of the pipe in conjunction with the spacing required between the pipes. Specifically, it will be difficult to install the '914 invention in vertical boreholes and trenches because the pipes will tend to become disarranged from their designed arrangements. This is especially true when the installation takes place in a vertical borehole filled with water. The '914 inventor suggests a solution. He uses spacers installed at intervals on the pipe. However, this increases the cost of assembly and transportation of the '914 system.\nAlso, the '914 invention uses an insulated pipe. Said insulated pipe does not contribute to the heat transfer process while occupying space in the borehole, and system efficiency is compromised.\nVertical borehole installations of the ground loop are usually required to be at least partially grouted. To optimize heat transfer it is common to grout the entire bore. This is accomplished by the insertion of a grout pipe all the way to the bottom of the bore. This grout pipe can sometimes be very difficult to insert into the bore as it has a tendency to catch on various irregular surfaces. The added spacers in the '914 configuration in conjunction with the space between pipes will make inserting this grout pipe cumbersome and laborious because the pipe may get caught on the spacers.\nApplicant's invention overcomes the problems with the art.\nSpecifically, Applicant twists 2 or more pipes together, making more efficient use of the bore hole space by placing more pipe into the borehole; thus creating more heat transfer surface area than the presently available pipe configurations. Because the 2 or more pipes are twisted together, said twisted pipes will support one another during installation. Further, said twisted pipe arrangement is easily rolled up and transported to a job site.\nAlso, the conduit of the present invention allows for greater thermal transfer from pipe to ground. If desired, the grouting material may be replaced by air or any insulating gas; thus, promoting pipe-to-pipe insulation.\nFinally, the conduit of the present invention can be manufactured from existing components and use in the market can begin almost immediately."}
{"text": "This application is based on the Japanese Patent Application 2000-054551, filed on Feb. 29, 2000, all the content of which is incorporated herein by reference.\n1. Field of the Invention\nThe present invention relates to a solid-state image pickup device and a method of manufacturing the same, and in particular, to a solid-state image pickup device and a method of manufacturing the same in which a solid-state image pickup element is attached to a package with high precision.\n2. Description of the Related Art\nWith an increasing need of higher picture quality for a digital video camera of integrated type and a digital still camera, research and development of a solid-state image pickup element for use with a solid-state image pickup device have been carried out for fine structure and high integration thereof. A need for an optical system is also highly increased, namely, an image must exactly focused on an image-pickup plane of the solid-state image pickup element.\nFIG. 8 shows a general configuration of a solid-state image pickup device in a cross section.\nIn FIG. 8, a solid-state image pickup device X includes a package 101 in a contour of a box and a solid-state image pickup element 103 housed in the package 101.\nAs can be seen from FIG. 8, the package 101 includes a bottom plate 101a and a sidewall section 101b rising from a peripheral section of the bottom plate 101a. The configuration defines an opening section 101c formed in an upper section of the package 101 and a housing section 101d surrounded by the bottom plate 101a and the side wall section 101b to house the image pickup element 103.\nThe image pickup element 103 is installed in the housing section 101d. The image pickup element 103 is connected via bonding wires 107 to lead wires of the package 101.\nThe opening section 101c of the package 101 is covered with a transparent shield plate 111 placed on an upper surface 101e of the side wall section 101b. \nThe solid-state image pickup device 103 ha a rear surface fixedly attached to the housing section 101d by adhesive 115.\nFIG. 9 is a cross-sectional view to explain a method of manufacturing a solid-state image pickup device.\nAs shown in FIG. 9, a package 101 on a carrier 123 is transported (to a predetermined position). The carrier 123 has an opening below the package 101.\nA package stage 121 made of iron-based material is moved upward through the opening to mount the package 101 on the stage 121. The stage 121 has a stage adsorbing hole 121a to fix the package 101 on the stage 121. The hole 121a produces a state like vacuum by sucking air to adsorb the rear surface of the package 101 on the package stage 121.\nThe housing section 101d of the package 101 is coated with adhesive 115. In the state, the solid-state image pickup element 103 is supported by a collet 125 and is placed at a predetermined position in the housing section 101d. \nThe collet 125 has an adsorbing hole 125a which adsorbs the solid-state image pickup element 103 to fixedly support the element 103. The collet 125 with the element 103 is moved downward to be pressed against the adhesive 115 to install the solid-state image pickup element 103 on the housing section 101d of the package 101. The adhesive 115 between the element 103 and the package 101 is cured or hardened by heat to fix the element 103 in the package 101.\nAfter a wire bonding process is conducted for the objective item, the upper opening section of the package 101 is covered by a shield plate to completely manufacture a solid-state image pickup device.\nFIG. 10 shows in a schematic cross-sectional view a configuration of an image pickup device including a solid-state image pickup device and an optical system including a lens.\nAs shown in FIG. 10, an image pickup device Y includes a solid-state image pickup device X, a solid-state image pickup device fixing member 131, an image pickup lens body tube 133, and an image pickup lens 135.\nThe fixing member 131 supports a bottom surface of the package 101. The lens body tube 133 in a cylindrical shape is formed on a peripheral edge section of the fixing member 131. On an inner circumferential section 133a of the body tube 133, an upper engaging section 133b and a lower engaging section 133c are formed to support the image pickup lens 135. The upper and lower engaging section 133b and 133c fixedly support the lens 135.\nAs the adhesive to fix the solid-state image pickup element on the package, there has been used thermosetting adhesive such as silver paste or resin-based adhesive. The silver paste has a high reflection factor for light and is therefore not suitable as adhesive for a solid-state image pickup element. Consequently, resin-based adhesive has been increasingly used.\nHowever, in a solid-state image pickup device including photo diodes in a fine structure, there arises a problem that resolution of the solid-state image pickup device deviates or fluctuates within a chip area of a solid-state image pickup device.\nTo increase picture quality of a digital video camera of integrated type and a digital still camera, only miniaturizing structure of photo diodes and the like is insufficient. To improve picture quality of a camera or the like, it is essential how to exactly adjust focus of an image pickup lens of a camera onto a light receiving plane of each photo diode formed in fine structure.\nFor this purpose, it is required to improve precision of tilt of image in a light receiving plane (parallelism of the light receiving plane with respect to the optical system).\nIt is therefore an object of the present invention to provide a technique to reduce fluctuation of the tilt angle of image in a solid-state image pickup element to increase picture quality of a camera or the like using the element.\nAccording to one aspect of the present invention, there is provided a method of manufacturing a solid-state image pickup device comprising the steps of preparing a package including a housing section to house a solid-state image pickup element chip and an opening section in an upper section thereof, sporadically applying adhesive with a predetermined thickness on a bottom surface of the housing section, moving the solid-state image pickup element toward the housing section, an upper surface of the package and an upper surface of the element each having desired parallelism with respect to a predetermined reference surface with high precision, bringing a rear surface of the solid-state image pickup element into contact with the adhesive and stopping the movement of the element before the element comes into contact with the bottom surface of the housing section, and curing the adhesive while the solid-state image pickup element is floating on the adhesive and fixing the element in the housing section.\nAccording to the present invention, there is provided a method of manufacturing a solid-state image pickup device, comprising the steps of preparing a package including a bottom plate and a side wall rising from the bottom plate to house a solid-state image pickup element chip, sporadically applying adhesive with a predetermined thickness on the bottom plate of the package, setting an upper surface of the side wall or an upper surface of the bottom plate as a reference plane, bringing a rear surface of the solid-state image pickup element into contact with the adhesive with parallelism adjusted with high precision between the upper surface of the element and the reference plane, moving the solid-state image pickup element, the rear surface of the element depressing the adhesive, the rear surface of the element not coming into contact with the upper surface of the bottom plate and curing the adhesive and fixing the solid-state image pickup element in the package.\nAccording to the present invention, there is provided a package comprising a bottom plate, a side wall rising from said bottom plate, a housing section for hosing a solid-state image pickup element chip, an opening section in an upper section of the package, and a projection projecting outward from said side wall, said projection including a flat upper surface.\nAccording to the present invention, there is provided a solid-state image pickup device comprising a package including a bottom plate, a side wall rising from said bottom plate, a housing section for hosing a solid-state image pickup element chip, and an opening section in an upper section thereof; adhesive with a predetermined width sporadically applied on said bottom plate, a solid-state image pickup element including a rear surface supported on the adhesive and being fixed thereon, said element being floating over a surface of said bottom plate; and a transparent shield plate for closing said opening section.\nIn the solid-state image pickup element, the fluctuation of resolution in the light receiving plane is reduced. This increases the picture quality of image pickup devices such as a camera using the element."}
{"text": "Mobile devices include cameras with relatively simple optics containing few camera elements. However, mobile devices also include increasingly powerful processors that can take advantage of computational photography techniques. For example, using computational photography, a series of photographs can be taken and processed to generate digital panoramas, high-dynamic-range (HDR) images, ultra-low-light images, and three dimensional “light field” images, among other types of images and effects.\nThe same numbers are used throughout the disclosure and the figures to reference like components and features. Numbers in the 100 series refer to features originally found in FIG. 1; numbers in the 200 series refer to features originally found in FIG. 2; and so on."}
{"text": "There is a need for personal communication systems that allow portable devices, such as pagers, telephones, computer terminals, and so forth, carried on the person to communicate with each other and with fixed location devices. For example, a user may want to store in a personal computer a message received over the air by a paging terminal. Known prior personal communication systems typically require that these two devices be interconnected by wire, which makes it cumbersome to affix them to users and/or interconnect them to each other, and thus, inconvenient to use.\nIn medical environments, systems for gathering information such as blood pressure, EKG readings, and so forth typically require that instruments taking the readings from a patient be connected, by wire, to a patient-carried system component that monitors or stores the information. These system are also cumbersome to affix to a user.\nIn other applications, wireless systems are currently used to transmit information between system components by, for example, radio waves, microwaves, infra-red signals and so forth. These systems may not be suitable for sending information between the user-mounted system components discussed above because of problems with interference in the immediate environment or between the signals transmitted from the various devices.\nFor example, devices in systems that use infra-red signals should optimally communicate with line-of-sight transmissions, which are not always possible between devices carried by a user. Further, the infra-red systems suffer from interference with ambient light, which can not always be controlled by the user. And, for systems transmitting signals at high frequencies, the users' bodies absorb the radiating energy, and thus, degrade the signals.\nFurther, such systems are subject to government regulation, since their signals radiate significantly. Also, these systems allow others to eavesdrop on the transmissions.\nWireless transmission systems have also been used to determine relative position. Such systems determine the position of a transmitter based on the timing or strength of signals received by various receivers. These systems are not well suited for and can be unreliable for determining position and orientation at close distances."}
{"text": "As shown in FIG. 30, output information (RGB image gray gradation image, CMYK image data) from a personal computer (PC) is supplied to the image processor provided in or connected with the image forming apparatus. In the image forming station, the output information from the PC is processed by a gamma correction portion for making correction in accordance with the gamma condition set by the printer driver depending on the user preferable, a converting portion to L*a*b* (only L* for gray), a converter to), C′M′Y′K′.\nA color space conversion from a many dimension to a many dimension, for example, from RGB to L*a*b* or the like, is effected by the image processor using a multi-dimensional conversion table information of an ICC profile.\nReferring to FIG. 31, there is shown a color setting portion of a printer driver which is a user interface (UI) provided for the user set a desired tone gradation. On the basis of the set point information, the gamma correction portion effects a gamma conversion using a unidimensional LUT.\nReferring to FIG. 30, the image data converted to the C′M′Y′K′ information is subjected further to a unidimensional LUT process by the engine tone gradation corrector. The engine tone gradation corrector is effective to keep the constant tone gradient of the printer engine. The state of the engine is detected using a patch image outputted on the recording material, that is, the sheet, and a unidimensional LUT is produced so as to provide a predetermined tone gradation curve. A stabilization control in which a patch outputted on the sheet thereafter is called automatic tone gradation correction.\nA pseudo-half-tone processing portion functions to reproduce a half-tone using a dither method or the like. Since the engine tone gradation property is different depending on the kind (line number and/or dot growth method) of the half-tone, the tone gradation corrector is provided with a LUT for each of half-tone processes.\nA technique assuring the stabilization property of color using the engine tone gradation corrector has been proposed.\nJapanese Laid-open Patent Application Hei 1-309082, for example, a patch pattern for toner content detection formed on the photosensitive member is read by a density sensor, and the read information is fed-back to the toner content controller in the developing device to provide a proper density.\nGenerally, such a toner patch is easy to produce and erase, but what is obtained thereby is only density information prior to image fixing, and therefore, when the control based on the toner patch is carried out, the influence of the secondary transfer or the fixing step are not reflected.\nJapanese Laid-open Patent Application Hei 1-309082 proposes that an image is read by a reader portion of a copying machine main assembly, and on the basis of the result of the reading, the image formation is controlled. With this method, the output print provided by fixing the toner on the sheet is read by the reader portion.\nAs a result of printing on standard, non-standard paper A (smooth thin paper), non-standard paper B (rough thick paper) under constant image forming conditions (charging, latent image and developing conditions), the density values detected by the reader portion are as shown in FIG. 32. On the basis of this information, a LUT shown in FIG. 33 is produced to provide a density.\nThe read density value is different depending on the paper surface of property, whiteness degree, thickness or the like, even if the toner amount thereon is the same, and therefore, the LUT varies. The LUT prepared on the basis of the non-standard paper B is lower in the maximum density than the target, and therefore, after input value of 240, data of 255 is used, for all. Therefore, over 240th level, the signal value exhibits no tone gradation. This is critical defect for the user who takes the tone gradient in the shadow portion in such a range important.\nIn addition, when the LUT is prepared on the basis of the non-standard paper A, the signal of 255 is outputted around the level of 240. Therefore, in the case of a signal indicative of solid image (halftone dot area percentage is 100%), a half-tone process pattern appears with the result that quality deterioration, and/or jaggy patterns (screen patterns) appear at edge portions. Thus, the quality of letters and/or lines deteriorates.\nIn order to solve such a problem, Japanese Laid-open Patent Application 2006-165752 proposes that when the input signal is 255 (halftone dot area percentage is 100), the OUT side of the LUT is made 255, thus providing a certain degree of tone gradation.\nFIG. 34 shows a LUT prepared on the basis of the non-standard paper B, and a LUT prepared by the method disclosed in Japanese Laid-open Patent Application 2006-165752. It will be understood that the LUT exhibits a smooth tone gradation. FIG. 35 shows a reproducibility of the density using the LUT. In the shadow portion, the density is different from the target density, but the tone gradient is slightly improved.\nHowever, with such technique, the tone gradation change of the shadow portion is so intense that an inflection point appears at (a) in FIG. 35, and therefore, smooth tone gradation property is not provided. Because of this, some professional users are not satisfied.\nAs another technique relating to the automatic tone gradation correction, Japanese Laid-open Patent Application 2004-289200 and Japanese Laid-open Patent Application 2006-222804 discloses that a patch image is printed on the non-standard paper using tone gradation correction data prepared on the basis of the standard paper, and the result of reading the image is used for a correction target. The automatic tone gradation correcting method is proposed as being free of the difference in the density or surface property of the paper per se.\nIn many cases, the density in the foregoing automatic tone gradation corrections is not an absolute density but is a paper based density (Null Density, relative density). In the following, the paper based density is called relative density. Relative density=absolute density of image−absolute density.\nThe density management method in the relative density is stipulated in ISO13656 or the like, and is an ordinary reference in the printing field Using the relative density, it is unnecessary to change the target density even when the printing is carried out on various kinds of paper, and therefore, the relative density has been used for management of the printing machines in printing companies\nAlso in the automatic tone gradation correction for an image forming apparatus, the automatic tone gradation correction has been carried out on the basis of the paper based density in view of variations of the reading apparatus and change with time thereof, differences of the paper density due to the differences of the paper lot difference, or the like\nHowever, in the case of using the automatic tone gradation correction disclosed in Japanese Laid-open Patent Application 2004-289200 and Japanese Laid-open Patent Application 2006-222804, the user has to prepare standard paper beforehand. In addition, the user has to be aware of the type of the non-standard paper and has to always carry out the tone gradation correction for the non-standard paper this means that what results is mere replacement of the standard paper with the non-standard paper. There is a liability that the non-standard paper registered on the basis of the standard paper may become out of stock. In such a case, the problem arises, because the user encounters the difference in the surface property, the whiteness degree, thickness, and/or grain of the paper due to lot difference, in addition, the difference in the whiteness degree, the paper per se due to the atmospheric corrosion resistance.\nWhen a plurality of kinds of paper are registered against the inventory depletion of the non-standard paper, the user has to make choice the kind of paper to use for the automatic tone gradation correction. With increase of the number of the prepared kinds, the operation becomes more difficult, and the probability of mistake increases. In addition, the memory capacity for the registration has to be large correspondingly with the result of the increase in cost\nFurthermore, the user is required to carry out twice the reading operations each including producing tone gradation correction data on the basis of reading of the printed patch on the standard paper, printing a patch on non-standard paper using the produced tone gradation correction data, reading the patch by the image reading station, and making correction. Therefore, the usability is not sufficient.\nMoreover, the coloring material of electrophotographic image is low in the transmission factor, and therefore, a solid image density is substantially immune to paper density. The conventional automatic tone gradation correction utilizing the density detection method in the printing field as it is, when the same amounts of the coloring material are on different sheets of paper, the relative densities are different. For this reason, when the automatic tone gradation correction is effected with different paper density, the result of the density after correction is different, and therefore, the automatic tone gradation correction with the standard paper is required."}
{"text": "1. Field of the Invention\nThis invention relates to a method and apparatus for the separation of a catalyst phase from a gas suspension phase in a fluidized catalytic cracking (FCC) unit. More particularly, it relates to an improved method and apparatus for separating the catalyst phase from the gas suspension phase in a closed cyclone process and system, to reduce the cyclone size requirements while maintaining the desired hydrocarbon and catalyst feed flowrate.\n2. Discussion of the Prior Art\nThe field of catalytic cracking, particularly fluid catalytic cracking, has undergone significant development improvements due primarily to advances in catalyst technology and product distribution obtained thereform. With the advent of high activity catalysts, and particularly crystalline zeolite cracking catalysts, new areas of operating technology have been encountered requiring refinements in processing techniques to take advantage of the high catalyst activity, selectivity and operating sensitivity.\nBy way of background, the hydrocarbon conversion catalyst usually employed in an FCC installation is preferably a high activity crystalline zeolite catalyst of a fluidizable particle size. The catalyst is transferred in suspended or dispersed phase condition with a hydrocarbon feed generally upwardly through one or more riser conversion zones (FCC cracking zones), providing a hydrocarbon residence time in each conversion zone in the range of 0.5 to about 10 seconds, and usually less than about 8 seconds. High temperature riser hydrocarbon conversions, occuring at temperatures of least 1000.degree. F. or higher and at 0.5 to 4 seconds hydrocarbon residence time in contact with the catalyst in the riser, are desirable for some operations before initiating separation of vaporous hydrocarbon product materials from the catalyst.\nRapid separation of catalyst from hydrocarbons discharged from a riser conversion zone is particularly desirable for restricting hydrocarbon conversion time. During the hydrocarbon conversion step, carbonaceous deposits accumulate on the catalyst particles and the particles entrain hydrocarbon vapors upon removal from the hydrocarbon conversion step. The entrained hydrocarbons are subjected to further contact with the catalyst until they are removed from the catalyst by a separator which could be a mechanical means, and/or stripping gas in a separate catalyst stripping zone. Hydrocarbon conversion products separated and materials stripped from the catalyst are combined and passed to a product fractionation step. Stripped catalyst containing deactivating amounts of carbonaceous material, hereinafter referred to as coke, is then passed to a catalyst regeneration operation.\nCyclonic equipment is typically used for efficient separation of fluidizable catalyst particles from the gas suspension stage. In U.S. patent application Ser. No. 612,277 filed May 21, 1984, entitled \"Closed Cyclone FCC Catalyst Separation Method and Apparatus\" by J. H. Haddad et al, a cyclone separation system is disclosed which reduces the undesirable extended residence time of product vapor within a large reactor vessel. The reduction in residence time is brought about by connecting a number of cyclone seperators in series directly to the output of the riser conversion zone with the gaseous output of the last cyclone conducting cracked hydrocarbons out of the reactor vessel. While extremely effective, the passage of virtually all of the hydrocarbon feed through the primary cyclone requires that it be designed for a relatively high volume flow of gas products as well as fluidized catalyst particles. This size requirement dictates that the reactor vessel, enclosing the riser conversion zone and separation cyclones, be relatively large.\nU.S. Pat. No. 4,404,095 to Haddad et al, discloses an FCC reactor comprising a riser 10 with radially extending sidearms 12 as the first catalyst-hydrocarbon separation means, as seen in FIG. 1 of the present application. The sidearms 12 force the suspension of the catalyst and hydrocarbons to change direction of flow, from the vertical to the horizontal direction, when striking end 14. The catalyst particles fall out of sidearm openings 16 in a downward direction, to a respective stripping means 18 which includes baffles 20. The hydrocarbons, with some entrained catalyst particles, proceed through passages 22 into the inlets 24 of cyclones 26. The sidearms 12 and the secondary separation means are enclosed by a vertical conduit 28 to prevent random uncontrolled thermal cracking of the hydrocarbons, which would ordinarly be present for a substantial period of time in the reactor vessel 30. Catalyst particles separated in cyclones 26 pass through diplegs 32 through to the bottom of the reactor and are removed through conduit 34. Hydrocarbon gases pass through the cyclone exhaust 36 into plenum 38 and out conduit 40 to a fractionation unit (not shown).\nIt can be seen in FIG. 1 that all of the catalyst particles and gaseous hydrocarbon materials exiting riser 10 must pass through the radially extending sidearms 12, requiring that they be sized to accommodate such a flow. Furthermore, during operation, substantial radial and longitudinal thermal expansion is encountered in the connecting conduits and intakes, and radially extending sidearms which necessitate some thermal accommodation."}
{"text": "The present invention relates to a system and method for providing communication among computers and devices in a network or distributed system, and more particularly it relates to a system and method of performing operations within the network or distributed system and transmitting throughout the network data content information.\nComputer users often desire information concerning changes in data through a computer network. In order to obtain updated information concerning particular updated data, however, a user must typically repeatedly access the information in order to obtain the most current version of it. For example, if the user desires to know when a particular stock price increases by five percent, the user typically must repeatedly access the information to obtain the current price, or obtain real time information concerning the stock price, and determine therefrom when an increase of five percent occurs. Likewise, if the user desires to sell the stock upon a five percent increase, the user typically must manually execute a sell command, rather than having the user\"\"s machine trigger execution of an automatic sell command upon detecting the required change. Also, the processing to determine when a particular change occurs typically happens at the user\"\"s machine, increasing the amount of processing required by that machine and other machines seeking the same information.\nObtaining information through the network also involves two-way communication. Each time a user desires updated information, the user\"\"s machine transmits a request for the information. In response, a server or another machine transmits the updated information to the user\"\"s machine. Therefore, receiving updates to particular information, such as stock prices, requires repeated two-way communication between user machines and a server, resulting in increased network traffic.\nIn addition, a computer may transmit information concerning particular data to an entire network of computers. For example, a server may transmit through a network information concerning current stock prices. However, it usually transmits that information throughout the entire network, making it available to any user who desires it. Such an approach increases network traffic in comparison to sending information only to users having requested it or with a specific need to know.\nAccordingly, a need exists for improved communication of changes in data within a computer network or distributed system."}
{"text": "1. Field of the Invention\nThe invention relates in general to a focusing controller and the method thereof, and more particularly to a focusing controller and the method thereof used in an optical disk drive.\n2. Description of the Related Art\nIn conventional operation, optical disc drives are subject to the control of servo signals. Among the servo signals of an optical disk drive, the focus error signal Fe is used to control the position of its pick up head, so as to focus the laser beam emitted by the pick up head on the optical disk. FIG. 1 is a block diagram of a convention optical disk drive. Conventional optical disk drive 100 includes a focusing controller 120 and a pick up head 110 to access an optical disk 130. The pick up head 110 outputs a focus error signal Fe according to the reflected light from the optical disk 130. The focusing controller 120 outputs a focus control signal Foo to control the pick up head 110 according to the focus error signal Fe.\nFIG. 2 is a diagram of the focusing controller, which mainly includes an analog-to-digital converter (A/D converter) 210, a focus compensator 220 and a digital-to-analog converter (D/A converter) 230. Functionally speaking, the A/D converter 210 is to receive and digitize the focus error signal Fe, for outputting a digital focus error signal to the focus compensator 220. Then, in subsequence, the D/A converter then converts the output of the focus compensator 220 to generate the focus control signal Foo for controlling the position of the pick up head 110.\nHowever, noises inside the optical disk drive affect the servo signals, including the focus error signal, so that the laser beam cannot precisely focus on the optical disk. There includes a plurality of tracks having pits and lands thereon used to record information on the surface of an optical disk. However, the focus error signal Fe is easily coupled with the track error signal Te during tracking operations, which indicates that the pickup head may not be allocated at an accurate position during the following focusing operations since the currently derived focus error signal Fe is erroneous because of coupling low-frequency error component from the track error signal Te. FIG. 3 shows waveform diagrams of the track error signal Te and the focus error signal Fe during tracking operations. Obviously, the focus error signal Fe is usually coupled with the low-frequency component of the track error signal Te such that the actually focus status of the pickup head may not obtained due to an erroneous Fe is derived."}
{"text": "1. Field of Invention\nThe invention relates generally to methods and apparatuses that process fluids or gases and, more specifically, to methods and apparatuses for extracting useful work from gases subsequent to combustion or to methods and apparatuses that move gasses or liquids.\n2. Art Background\nProcessing liquids and gases during the extraction of useful work, such as in an internal combustion engine or while transferring a liquid or a gas from point A to point B, has been the subject of much activity. For example, deriving useful work from the combustion of fuels has been of value to mankind for many years. Internal combustion engines are known in the art and are used in a wide range of devices, such as grass trimmers, automobiles, and large scale earth moving equipment such as trucks, bull dozers, etc.\nMany of the current internal combustion engines are based around a reciprocating piston design or a translating rotor. In a reciprocating piston design, a piston travels back and forth in a cylinder moving to and fro, reversing direction at each end of its range of motion within the cylinder. A piston has a finite mass and thus energy is required to change the direction of the piston within the cylinder as the piston travels to and fro. This energy is deducted from the energy developed during the operation of the engine and necessarily reduces the amount of energy available as useful work output of an engine. Energy directed to non-useful work output contributes to the inefficiency of an engine.\nInternal combustion engines typically burn fossil fuels and expel products of combustion. Fossil fuels exist in limited quantities and are viewed as non-renewable resources. Society is in general agreement that improved fuel efficiency is a useful and beneficial goal. Since fossil fuels exist in limited supply, the current fuel efficiency of such engines may present a problem. Additionally, the products of combustion are a source of environmental pollution and may also present a problem.\nInternal combustion engines also employ valves to interrupt the flow of gas during the various cycles of operation, such as intake and exhaust valves. Thus, gas accelerates and decelerates during the operation of the engine resulting in non-uniform flow. Such non-uniform flow can contribute to engine inefficiency via the inertial effects described above in conjunction with the piston or the translating rotor as well as other effects not discussed herein. Such non-uniform flow may present a problem.\nPumps are known and used in the art to move materials such as viscous materials, fluids, and/or gasses from one location to another. Pumps are made using the reciprocating piston design discussed above. Work input to a pump is required to operate the pump. Therefore, energy inefficiency can result from a reciprocating piston pump design. Similarly, diaphragm pumps employ a moving diaphragm which changes direction continuously, resulting in an expenditure of energy to change the direction of the diaphragm; such inefficiency may present a problem."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a surgical instrument and more specifically it relates to a surgical skin punch apparatus to excise subcutaneous abnormalities and allow for efficient resealing of the skin surface.\n2. Description of the Related Art\nAny discussion of the related art throughout the specification should in no way be considered as an admission that such related art is widely known or forms part of common general knowledge in the field.\nSubcutaneous abnormalities, such as lipomas, cysts, lesions, etc., that desire removal are often surgically removed using various types of tools. Often times, a “free hand” scalpel or a circular-blade punch is used to create a surgical incision in the skin about one-half the size of the tumor or other mass that is desired to be removed. The removal of the tumor through the small incision or punch decreases the size of the resulting scar, but does not eliminate it completely. In addition, when using the “free hand” approach the resulting scar can often times be uneven or appear deformed.\nWhen using a circular-blade punch, one method is to discard of the circular skin fragment and leave the site to heal. However, this method generally results in a larger than necessary scar and may be disfavored by patients. Another method is to replace the circular skin fragment back over the defect.\nHowever, limitations of this technique include the possibility that the circular skin fragment may not heal flush with the surrounding skin and thus result in a depressed or elevated scar, or the possibility that blood will not adequately flow to the replaced circular skin fragment which might lead to necrosis and scarring. Because of the inherent problems with the related art, there is a need for a new and improved surgical skin punch apparatus to excise subcutaneous abnormalities and allow for efficient resealing of the skin surface."}
{"text": "As is known in the art, Bluetooth is a wireless technology that allows short range communications between Bluetooth enabled computing devices. Bluetooth enabled computing devices can include portable and/or fixed computing devices.\nAccordingly, Bluetooth provides short range radio communications between computing devices located within a predefined physical range. Typically, the short range of Bluetooth technology is application specific and can be pre-defined. Bluetooth operates in the 2.4 GHz ISM radio frequency band. Typically, Bluetooth computing devices operate in a master-slave configuration such that a master Bluetooth computing device can communicate with a number of slave computing devices (i.e. within a piconet or Bluetooth communication network of devices). In turn each of the slave computing devices can act as a master within a separate piconet and connect to a number of designated devices.\nConventionally, Bluetooth security measures have been limited to a few main modes of security. In one mode, no authentication means is defined for the computing device such that anyone can access the device. In another mode, a user can switch the Bluetooth connection mode to “non-discoverable” and avoid being discovered by other Bluetooth computing devices entirely. In this mode, the device could also be set such that the non-discoverable mode occurs for a predefined period of time. In yet another mode, users of different Bluetooth computing devices may need to authenticate by exchanging a password or PIN number between one another after which a connection can be established.\nHowever, all these types of security are limited in their capabilities and in some cases could still allow unwanted connections to be made. Also, it may not be desirable for a computing device to define themselves as non-discoverable as this would limit the connections that occur subsequently. Accordingly, it would be advantageous to provide an improved method of controlling access between a computing device and other computing devices in a Bluetooth communication network.\nAccordingly, it is desirable to provide a solution that obviates or mitigates at least some of the above mentioned disadvantages."}
{"text": "The present invention relates to a universal mechanical control system for moving a controlled device in response to the movement of a universally movable control rod. The invention is particularly applicable for controlling a laser beam surgical scalpel in accordance with the movement of a manipulatable joy stick and is therefore described below with respect to this application.\nLaser beam surgical scalpels are increasingly being used for performing delicate surgical operations. A known type of such scalpels includes a microscope to permit viewing by the surgeon of the working area, and an optical control unit including a control rod, commonly called a joy stick, manipulatable by the surgeon for directing the laser beam with respect to the working area. Critical factors in the successful use of such scalpels for performing delicate surgery are the ease and precision by which the laser beam may be manipulated by the surgeon by using the joy stick. Various types of mechanical and electrical control systems have been proposed or are now in use for this purpose, but the known systems are generally of complicated, bulky, and expensive construction and/or are not entirely satisfactory with respect to the ease and precision by which they may be conveniently manipulated by the surgeon during the performance of a delicate surgical procedure.\nAn object of the present invention is to provide a simple, inexpensive and easily-operated mechanical control system for moving a controlled device in response to the movement of a universally movable control member. Another object of the invention is to provide a laser beam surgical scalpel with the above mechanical control system for manipulating the laser beam with respect to the working area in response to the manipulations of a joy stick."}
{"text": "This invention relates to an FM stereophonic receiver having an ARC (automatic reception control) function.\nFIG. 7 is a block diagram showing the arrangement of essential parts of a conventional FM stereophonic receiver.\nIn FIG. 7, reference numeral 1 designates an FM detector circuit which detects an intermediate frequency signal IF applied thereto, to output a main signal S.sub.M and a sub signal S.sub.S.\nReference numeral 2 designates a high cut circuit which operates as follows: The high cut circuit 2 controls the frequency characteristic of the main signal S.sub.M outputted by the FM detector circuit 1 in accordance with a high cut control signal S.sub.H provided by a high cut control circuit 8. The high cut control circuit 8 outputs control signal S.sub.H according to a received signal field strength level. The high cut circuit 2 outputs a high cut signal S.sub.H1 according to the main signal S.sub.M and the control signal S.sub.H.\nReference numeral 3 designates a variable attenuation circuit. The variable attenuation circuit 3 attenuates the sub signal S.sub.S outputted by the FM detector circuit 1 in accordance with a separation control signal S.sub.SP provided by a separation control circuit 3A. The separation circuit 3A outputs control signal S.sub.SP according to a received signal field strength level. The variable attenuation circuit 3 provides an attenuation signal S.sub.A in accordance with the sub signal S.sub.S and the control signal S.sub.SP.\nReference numeral 4 designates a stereophonic demodulation circuit which operates as follows: The stereophonic demodulation circuit 4 separates right and left channel signals S.sub.R and S.sub.LA according to the high cut signal S.sub.H1 outputted by the high cut circuit 2 and the attenuation signal S.sub.A outputted by the variable attenuation circuit 3, and outputs the right and left channels signals thus separated.\nThe FM stereophonic receiver thus organized operates as follows:\nThe separation control circuit 3A outputs the separation control signal S.sub.SP according to the received signal field strength level in such a manner that, when the received signal field strength level is high, the amount of attenuation of the variable attenuation circuit 3 is decreased, whereas when the received signal field strength level is low, the amount of attenuation is increased. Thus, the amount of attenuation of the variable attenuation circuit 3 continuously changes in accordance with the received signal field strength level.\nTherefore, the sub signal S.sub.S outputted by the FM detector circuit 1 is attenuated according to the received signal field strength level, and the right and left channel signals outputted by the stereophonic demodulation circuit 4 is improved in S/N ratio.\nThe separation control is effected when the received signal field strength level is within the range of, for instance, from 45 dB/.mu.V to 25 dB/.mu.V, whereas the frequency control is carried out when the received signal field strength level is less than 35 dB/.mu.V.\nThe conventional FM stereophonic receiver designed as described above suffers from the following difficulty: Since the sub signal S.sub.S is attenuated according to the received signal field strength, when the received signal field strength level becomes relatively low, the sub signal S.sub.S is completely attenuated, with the result that the stereophonic broadcast turns into a monaural broadcast."}
{"text": "This invention relates to power amplifiers and, more specifically, to a circuit and method for adjusting the voltage at the commonly connected ballasted ends of a plurality of base ballasted HBTs forming a power amplifier when the voltage appearing at one of the bases of the plurality of HBTs drops below a threshold.\nHetero-junction bipolar transistor (HBT) power amplifiers are becoming the standard for cellular applications due to their high power density and reduction in die size. HBT\"\"s can draw substantial base current during high power operation. Since multi-finger power devices suffer from thermally related current collapse they cannot be operated without ballast resistors on either the emitter or base. Emitter ballasting is not practical for power amplifiers because of the high emitter currents and small resistor values, so base ballasting must be used. Because an individual HBT is a very small device, it is required to be paralleled with multiple HBTs to achieve high power operation required for most power amplifier applications.\nAccordingly, HBT power amplifiers such as those used in radio frequency (RF) applications employ multiple small devices connected in parallel. As noted above, each of these small HBT devices require a ballast resistor to be connected to its base before being connected to the other HBTs forming the power amplifier (PA). The resistive ballasting of individual cells keeps parallel HBT fingers from entering thermal collapse. Additionally, a capacitor may be used to bypass the base resistor to preserve high frequency gain or the RF signal may be fed to the base connections through a separate capacitor. For purposes of illustration, the figures contained herein will illustrate the principal using capacitor bypassed ballast resistors although those skilled in the art will realize that this embodiment of the invention will work the same regardless of the connection of the RF capacitors feeding the base connection.\nFIG. 1 shows a typical multi-fingered base ballasted Power Amplifier (PA) circuit. A plurality of HBTs 110, each ballasted with a resistor 130/capacitor 120 are connected in parallel. For each small HBT device, first ends 130a, 120a of a resistor 130 and a capacitor 120 are connected to the base 185 of the individual HBT devices and the other ends 130b, 120b of the resistor and the capacitor become the input 180 of each base ballasted HBT device 190. For purposes of simplicity, when two or more base ballasted HBT devices 190 are xe2x80x9cconnected in parallelxe2x80x9d, their collectors 160 share a first common node, their emitters 170 share a second common node connected to ground, and the inputs 180 share a third common node. A radio frequency signal is received at the input 140 and connected to the commonly connected inputs 180 of the base ballasted HBT devices 190. The commonly connected collectors 160 that are connected to a voltage source 155 produce an amplified RF output 150.\nDue to the base current requirements, a biasing circuit 195 is usually included. Typical biasing circuits with RF decoupling components neglected for simplicity are shown in FIGS. 2 and 3. FIG. 2 shows a diode biasing circuit 200. The base of an HBT device 210 is connected to the collector and the collective inputs 180 of the of the base ballasted HBT devices 190 of FIG. 1. The emitter is connected to ground. The first end 220b of a reference resistor 220 is connected to the collector and base while the second end 220a of the reference resistor 220 is connected to a reference voltage 230.\nFIG. 3 shows a preferred current mirror biasing circuit 300. A current mirror is formed by HBT devices 310 and 320. The collector of the first HBT device 310 is connected to a voltage source 350, its emitter is connected to the base of the second HBT device 320 and its base is connected to the collector of the second HBT device 320. The emitter of the second device is connected to ground. And finally, the first end 330b of a reference resistor 330 is connected to the base of the first HBT device 310 and the collector of the second HBT device 320 while the second end 330a of the reference resistor is connected to a reference voltage 340.\nAlthough not exhaustive, these biasing circuits are typical of those employed in the industry although other types of biasing circuits are contemplated and may be used with the invention. The biasing circuits try to keep the current through the power device constant with variations in temperature and reference voltages. Although either of these biasing circuits or others could be used, current mirror biasing is typically preferred and will be used in the discussion.\nUsing the current mirror of FIG. 3 as the Bias of FIG. 1, the reference voltage 340 and the reference resistor 330 form a constant current source which is mirrored by the first HBT device 310 and the second HBT device 320. If no ballast resistors 130 were required, the current mirror would be adequate up to the limits of the HBT devices 310 and 320. However, with ballast resistors and during high power operation, the current mirror is unable to keep the voltage on the bases 185 of the individual HBT power device cells 110 constant because of the drop on the ballast resistor. In power operation when more HBTs are connected in parallel, increased base current is required from the current mirror. This strain on the current mirror results in increased voltage drops across the ballasting resistors 130 resulting in the voltage at the base of the individual HBT devices 110 to droop, limiting linearity and maximum output power.\nWhat is required is an improved HBT power amplifier circuit that doesn\"\"t effect the quiescent point at lower output powers, but comes into play when higher powers are being generated that effectively prevents this drooping from occurring."}
{"text": "1. Field of the Invention\nThe present invention relates to a head-mounted image display apparatus.\n2. Description of the Prior Art\nConventional methods of obtaining high realism in an image display apparatus that directs image light based on a displayed image to the user's pupil include a method in which a displayed image is projected onto the user's pupil as a virtual image without an intermediate image being formed and a method in which a displayed image is projected onto a screen and the user views the image. These methods are advantageous in that a wide viewing angle and a large pupil diameter are comparatively easily obtained and high realism is obtained.\nHowever, in the method in which a virtual image is projected onto the user's eye without an intermediate image being formed, a multiplicity of large-diameter lens elements are necessary to achieve a wide viewing angle and a large pupil diameter, which increases the cost of the optical system and the weight. On the other hand, the method in which an image is projected onto a screen and the user views the image is advantageous in that a large pupil diameter can be achieved by diffusion by the screen and that the projection optical system can be made compact. FIG. 7 schematically shows the structure of the optical system of a conventional head-mounted image display apparatus based on this method.\nAn image display element 53 is illuminated by a lighting source 52, and the image light based on the image display exits from a projection optical system 54. The image light is reflected at a half mirror 51 and projected onto a reflective screen 56 to form an image. The reflected light from the screen 56 is transmitted by the half mirror 51 and supplied to the user's pupil EP. By this light, the user views the image on the screen 56.\nHowever, in the conventional image display apparatus using a screen, the screen is fixed in front of the user. That is, since the screen is disposed away from the head, the overall mass moment of the apparatus imposed on the user's head is great, so that the load on the user when he or she wears the apparatus on his or her head is heavy."}
{"text": "Different types of telecommunications networks are known to the art. The typical telephone network offers telephony services. The telex network offers the telex communications service. A circuit switched data network offers data communications services, the paging network offers the paging service. The typical telecommunications network can be used as a carrier for other services, such as telex and data transmission, for instance. The packet switched network may be a carrier of different services, such as Datapack, Teletex, Videotex and Datel. Traffic in one network can be adapted for communication to destinations in another network, by introducing gateways or ports in the first-mentioned network. ISDN (Integrated Services Digital Network) is an integrated service network which is able to offer several of the aforesaid services. The subscribers connect all of their equipment to the same \"wire pair\" in ISDN.\nA user will traditionally use different equipment, herein-after called terminals, and different procedures to obtain access to different networks. The function of the terminal is to make the services of the network available to the user. A given terminal has the functionality for interaction with a given network, the functionality for interaction with a user, and the functionality of coordinating the behaviour of the terminal to the user and the behavior of the terminal to the network. User interaction with the terminal is different and will depend on the network to which the terminal has access. This means that the user must be aware of the access procedure for each of those networks with which he wishes to communicate through his different terminals. This means in practice that one access procedure is used when telephoning, that another access procedure is used when using a mobile telephone, and that a third access procedure is used when wishing to page a subscriber. The supplementary services of a telephone network can be utilized by dialling or keying a specific key sequence. The same supplementary service, for instance unconditional call forwarding or conditional call forwarding, have different key sequences in different networks.\nA common feature of all telecommunications networks is that they enable a connection to be established between two subscribers. Different networks are able to set up a connection between two terminals in the different networks, even when one network A cooperates with another network B through a port. On the other hand, none of the telecommunications networks is aware of why the connection is established between two terminals. Such knowledge lies beyond the makeup of the networks. Thus, if a subscriber X wishes to communicate with a subscriber Y and subscriber Y is located in a third network C which is different from A and B, none of the networks A or B is able to initiate the setup of a connection to subscriber Y in the network C.\nA user Y who can be accessed in several different networks has a number of possibilities from which to choose in present-day techniques. Most typically the user Y does nothing, implying that a calling party X will have to search the called subscriber Y by making a number of calls in each of those networks that subscriber Y believes that subscriber X can be found. Another possibility for the user Y is to control incoming traffic to his/her/its present destination, by activating some form of redirection or diversion in one or in a number of the networks. In the following redirection and diversion are used as synonymous terms."}
{"text": "1. Field of the invention:\nThis invention relates in general to flexible pipes for conveying fluids at high pressure, and in particular to those used for subsea oil and gas drilling and production.\n2. Description of the Prior Art:\nIn subsea well production and drilling operations, flexible pipes are used for various purposes such as underwater flowlines. A flexible pipe is herein considered to be a pipe having, at least during installation, sufficient flexibility along its longitudinal axis to accept a minimum radius of curvature in the order of five to twenty times the diameter of the pipe, without weakening the structure.\nThese pipes must resist combined hoop and longitudinal end loads from the high pressure conveyed fluids as well tensile loading during installation and operation. These pipes must also withstand crushing due to external pressures with or without internal pressure inside. Once in place, the internal pressure makes the pipe stiffer, and flexibility is no longer needed. A flexible pipe is easier to transport and install than a rigid pipe.\nA typical flexible pipe structure comprises at least one fluid containing tube made of polymer material. A reinforcing armor of metal wires is wound spirally around the fluid containing tube. In some instances, a third structure consisting of a helical winding of interlocking strips is located inside the containing tube. This last structure prevents the fluid containing tube from collapsing under flexible pipe annular pressure and provides radial support to the reinforcing armor while the flexible pipe is under tensile loads without internal pressure. Improvements in strength, manufacturing ease, and cost are desirable."}
{"text": "1. Field of the Invention\nThe invention relates to cycloidal drives, and more particularly to drives for reciprocating parts which require speed reduction during intermediate portions of their strokes. An example of such driven parts is a mechanical welding press which requires a mid-stroke slow down so as to minimize shock to parts when they are being lifted from automation rails, and reduce the noise level inherent with this shock condition. Another example is an automation system in which a reciprocating shuttle on a conveyor must be slowed down at mid-stroke for similar purposes.\n2. Description of the Prior Art\nGervais U.S. Pat. No. 3,559,767 shows a planetary gear system with a crank arrangement describing an epicycloidal path in order to reduce the platen speed at certain points in the motion. This patented construction has a number of disadvantages as compared with the apparatus of this invention. It utilizes two fixed external sun gears and pin-ended joints, thus requiring two primary gear assemblies and making it impossible, because of the geometry of the construction, to support the bull gear shafts at opposite ends. Instead they must be cantilever mounted. It is further believed that the Gervais construction would not be interchangeable with most present day welding press configurations.\nMatsushita U.S. Pat. No. 3,971,261 pertains to a slide driving device for use in metal working presses which is of the cycloidal type. This patent utilizes an external fixed gear and exhibits certain disadvantages as compared with the present invention which will become apparent from the following description.\nSeki U.S. Pat. No. 3,956,942 also illustrates a cycloidal arrangement in which there is a reduction of velocity in the center of the stroke. However the mechanism is substantially different than that of the present invention, and advantages of the present invention wll become apparent below.\nThe following patents are of interest insofar as they pertain to either cycloidal gear trains or to gear train arrangements having some features pertinent to the present invention:\nStaecker et al: U.S. Pat. No. 2,822,696 PA1 Fracke: U.S. Pat. No. 3,468,173 PA1 Koderman: U.S. Pat. No. 3,886,805 PA1 Hosoda et al: U.S. Pat. No. 3,921,465 PA1 Mori et al: U.S. Pat. No. 3,960,024 PA1 Obra: U.S. Pat. No. 4,020,708 PA1 Obra: U.S. Pat. No. 4,089,228 PA1 Drungil: U.S. Pat. No. 4,095,522"}
{"text": "Conventionally, there is known a vane pump which includes: a housing having an approximately circular pump room; a rotor which rotates at an eccentric position relative to the center of the pump room; and a vane rotated by the rotor for dividing the pump room full-time into a plurality of spaces.\nThen, in order to lubricate such a vane pump, there is known a vane pump configured in a manner that an oil supply passage which intermittently communicates with the pump room owing to rotation of the rotor described above is formed in the rotor, and a lubricating oil is intermittently fed from the oil supply passage into the pump room (Patent Document 1).\nHowever, in the case of the vane pump having such oil supply passage, when the rotor stops with the oil supply passage being in communication with the pump room, then owing to a negative pressure in the pump room, the lubricating oil in the oil supply passage is sucked down into the pump room, and when the vane pump, subsequently, gets started, the vane may be damaged by an excessive load which is applied to the vane to discharge this lubricating oil.\nTo address such a problem, there is known a technology that an air passage constantly communicating with an atmospheric air is formed in the oil supply passage, and when the rotor stops, a negative pressure in the pump room is eliminated by sucking an atmospheric air into the pump room through the air passage, thereby a large amount of the lubricating oil is prevented from flowing into the pump room (Patent Document 2).\nPatent Document 1: Japanese Patent No. 3107906 (particularly see paragraph 0022)\nPatent Document 2: Japanese Patent Laid-Open No. 2003-239882 (particularly see paragraph 0012)"}
{"text": "1. Field of the Invention\nThe present invention relates to a display apparatus and a display method. More particularly, the present invention relates to a display apparatus capable of displaying video providing a sense of perspective in accordance with the line of sight of a user and a display method.\n2. Description of the Related Art\nIn related art, enhancement processing adapted to video itself to be output to a monitor has been performed. Here, the enhancement processing is processing to improve sharpness of video so that sharper video is presented to the user. However, enhancement processing adapted to video itself to be output to a monitor does not have in most cases three-dimensional information added to the video itself and thus, the video after being processed may be different from what the user likes.\nMoreover, for example, in Japanese Patent Application Laid-Open No. 2000-221953, a technology to perform image processing in accordance with a detected result after shooting a user by a camera and detecting the line of sight of the user using the shot image is disclosed.    [Patent Document 1] Japanese Patent Application Laid-Open No. 2000-221953"}
{"text": "In recent years, the electric vehicles have achieved great improvements in the core technologies and have been gradually entering the stage of practical use and small scale industrialization.\nAs one of the key components of electric vehicle, the OBCM (on board charge module) is an electronic device mounted on board for charging battery packs. The OBCM uses input wires to receive alternative current from the AC power grid, and outputs high voltage direct current to charge the on board high voltage battery pack, and uses a communication port to maintain real-time interactive communications with a battery management system (BMS). Due to the severe on board operation conditions, the structure and packaging design and as well as the thermal management are the key challenges for OBCM design."}
{"text": "In recent years, automobile manufacturers, insurance companies, and others have begun equipping and/or retrofitting vehicles with a variety of sensors, interfaces, communicative connections, etc. for monitoring the operation of the vehicles. These systems can provide insights into how a vehicle is being driven, where a vehicle is being driven, if components of a vehicle are failing or need maintenance, and much more. However, many current monitoring systems are installed by the manufacturer of a vehicle or retroactively coupled to data ports, on-board computers, or other couplings built into a vehicle by the manufacturer. As such, vehicles that do not come equipped with monitoring systems from the manufacturer and/or do not possess certain types of data ports may not be able to utilize these monitoring systems.\nEven if monitoring systems can be coupled to data ports or on-board computers of a vehicle, these attached monitoring systems could have some impact on the performance of the vehicle, given that they are coupled to a data port or on-board computer of the vehicle. For example, an attached monitoring system, attached to a data port or on-board computer of a vehicle, may introduce unexpected current, voltage, or other signals into the electrical and computer system of the vehicle. These unexpected signals may cause the vehicle to stall or operate in a non-ideal manner."}
{"text": "Wi-Fi networks are widely-available and used more often with the growth of mobile devices. These networks operate under standards such as IEEE 802.11 (promulgated by the Institute of Electrical and Electronics Engineers), and others (e.g., Hot Spot 2.0). Management of these networks will become more burdensome and complex with increasing loads and numbers of connections.\nWi-Fi controllers are coupled to access points of a communication network in order to manage connections with stations seeking access to a backbone network, including the Internet. One aspect of management is the assignment and tracking of BSSIDs (Basic Service Set Identifiers) that identify access points to stations for traffic flows. SDN in the wired backbone supports the wireless portion of a communication network by centrally managing a data plane or layer 2 routing decisions.\nHowever, the benefit of centrally managing data plane behavior comes with a burden of tunneling packets from an access point to an SDN controller. The additional overhead can slow down network performance with additional packet transfers adding to network congestion and processing loads. The effect can be exasperated by flows of local devices sharing a common access point, such as Bonjour services transmitting to nearby stations.\nWhat is needed is a robust technique to intelligently bridge certain flows instead of tunneling."}
{"text": "The present invention generally relates to an apparatus and method for executing communication protocol conversions and, in particular, relates to such an apparatus and method whereby conversion from an external protocol to a uniform network protocol is efficiently executed.\nIn an effort to simplify what is frequently a rather complex design problem, most telecommunication systems are organized via the use of layers, or levels. Each layer is independently structured based on the services and functions provided by the adjacent layers. The number of layers in any given system or subsystem depends on the specific functions and services to be provided thereby. In addition, each layer is designed to communicate with adjacent layers via an interface. Generally, each layer provides certain clearly defined services or functions and, consequently, can be effectively designed fairly independent from the adjacent functions or services. Hence, the detailed implementation of the defined services and functions of any given layer are more simply designed. As one example, the first layer, or level, can be defined as the actual, i.e., physical, communication medium of the system, such as a coaxial cable, a printed circuit strip line, a microwave transmission network, or the like. In actually, most telecommunication systems include more than one of the above-mentioned communication media.\nEach layer is usually provided with a unique set of rules that regulate and control data communication within that layer and define the parameters for all layers interfacing therewith. In fact, any particular layer may include a plurality of nested sets of communication rules whereby, for example, a predetermined one of different data communication medium is selected according to the data destination and the appropriate set of rules are enforced. In any data communication device, whether a communication terminal or a worldwide communication network, the functions and services of each layer, as well as the transfer of data between adjacent layers are frequently regulated, i.e. defined, by the operating system software. By definition, the totality of the rules and regulations governing communication with such a device, regardless of the number of internal or nested layers provided therewithin, is known as the communication protocol of that device.\nHowever, in order for two devices, such as two different terminals, having different communication protocols to communicate with each other, a rather complex protocol conversion must be performed. Generally, in a communication network, this conversion can be executed in a number of different ways. In one conventional technique, a centralized processor having all necessary protocols stored therein, i.e. the communication protocol for each different type device, is included in the network. A centralized processor technique is frequently employed because of the magnitude of the memory required for storing all possible communication protocol conversion information for all protocols being used in any given network. In such an approach, data from a source device is first transferred to the centralized processor wherein the communication protocol conversion is executed. Thereafter, the data is transferred to the destination device having a communication protocol compatible therewith. However, as a consequence, all data entering such a network must pass through the centralized processor for executing the necessary communication protocol conversions and the transference of the data to the addressee. The data path through the central processor is thus traversed by all data entering the network, even if the transmitting and receiving peripherals have identical communication protocols. Understandably, this is inefficient, both because it is time consuming when no communication protocol conversion is necessary but also because the potential for data transfer errors increases as the network path length and transmit time therealong increases.\nThe most prevalent alternative scheme of executing communication protocol conversions between external peripheral devices operating under different protocols requires each subsystem of a network to execute the requisite communication protocol conversions. In such a scheme, each subsystem is provided with a dedicated subsystem processor containing all communication protocol conversion information for the protocols acceptable to the system. For example, if four different communication protocols are predesigned into a data communication network, each subsystem interconnected thereto must include a processor having the capacity, both in memory and actual processing capability, to convert between any combination of any two of those four protocols. Such conversions, of course, are in addition to the usual communication protocol conversion to the network protocol. In addition, similar to the central processing technique discussed above, all data entering, from either a peripheral or the network, a subsystem must pass through the subsystem processor. Thus, this technique is inefficient since the subsystem processor is a source of bottlenecking.\nIn addition to the above-mentioned drawbacks, the greatest and most expensive difficulty of either technique is the reprogramming required if a peripheral with yet another protocol is to be interconnected with the network. In fact, the expense of such reprogramming is often prohibitive, thus limiting such networks to interfacing with only specific types of peripherals. It is for just this reason that many data communication device manufactures emphasize the compatibility among their own products.\nFrom the above it is apparent that an appparatus and a method for executing communication protocol conversions between different communication protocols associated with external devices is highly desirable in light of the plethora of different data communication devices presently available and anticipated for providing data communication services to a burgeoning market."}
{"text": "Aesthetics of any vehicle is a one of the prime requirements and concerns for any automobile user ranging from the tires to the hood of that vehicle, and other parts as well. For example, the shine and glossy facade of the wheels enhances the overall splendor of an automobile. However, the luster of the wheels shine is dulled by various factors including brake dust, road grime, muddy water, etc. Out of those, brake dust not only deteriorates the glossiness, but also the dusts long term contact augments the corrosion tendencies of the wheel.\nGenerally, brake dust is a product of wear from the friction between the brake pads and rotors of the vehicle. Brake pads convert the kinetic energy of the car to thermal energy by friction. Two brake pads are contained in the brake caliper with their friction surfaces facing the rotor. When the brakes are hydraulically applied, the caliper clamps or squeezes the two pads together into the spinning rotor to slow/stop the vehicle. When a brake pad is heated by contact with a rotor, it transfers small amounts of friction material to the disc, turning it a dull gray. The brake pad and disc (both now with friction material), then “stick” to each other, providing the friction that stops the vehicle. In order to maintain the original appearance of automobiles, a wide variety of automobile appearance products are used in consumer applications. Specifically, numerous wheel and tire cleaners are commercially available to maintain the appearance of soiled tires, metal or plastic wheels, wheel covers, or hubcaps. These wheel cleaners are typically acidic formulations, with pH ranges from 2 to 4. The extreme acidity of the wheel cleaner is required to effectively remove brake dust, road soil, and grime. These formulations also typically use anionic, amphoteric, or nonionic surfactants to enhance the detergency of the wheel cleaning formulation. The extremely low pH of these formulations makes their compatibility poor with many metallic wheels, and most of these products are devised for use on clear-coated wheels so there is no (or very low) chance of acid corrosion. The compositions must be washed or wiped from the tires and wheels within 15 seconds to one minute to prevent permanent damage to the wheels, and also to facilitate cleaning.\nThe principle of self-cleaning plays a major role in brake dust repellent coatings. In order to achieve a good self-cleaning and dirt repellent surface, the surface not only has to be very hydrophobic, but also has to have some degree of roughness. A suitable combination of structure and hydrophobic properties does not permit any kind of dust or contaminants, including water, to adhere on the surface.\nA number of attempts have been made, as evidenced in the prior art, to obtain a super hydrophobic and antistatic composition. U.S. Pat. No. 7,083,828 to Muller et al. describes a formulation for detachable dirt and water repellent coating which has to be applied at room temperature only. U.S. Patent publication number US 2007/0190308 to Brand et al. discloses the formulation for scratch resistant and self-cleaning protective layer for rims of automobiles, but has the drawback that the formulation requires two layers to achieve the effect, and has to cure at an elevated temperature (up to 200° C.). U.S. patent publication number 2009/0018249 to Kanagasabapathy et al. discloses a coating composition for a hydrophobic self-cleaning effect but has the drawback that the formulation has to be applied on wax surface only. U.S. patent publication number 2004/0213904 to Muller et al. discloses a process for producing a water repellent and dirt detachable surface coating consisting of hydrophobic particles in conjunction with alkyd modified silicone wax in a highly volatile siloxane. The suspension is applied to the surface of an article, and volatile siloxane is evaporated forming a grayish coating. The disadvantage of this coating is hazy layer that is left on the treated surface, and is not suited for cosmetic applications.\nFurthermore, there are several published articles which have dealt with elimination of water marks. U.S. Pat. No. 5,759,980 to Russo et al. discloses a coating comprised of silicone based surfactant and a polymer, where the Polymer is used as a binder. However, this coating is hydrophilic in nature and may tend to be removed from the surface by single wipe.\nThe patent publication number WO 2008/153687 discloses a super hydrophobic coating composition composed of nano silica and an acetone solvent combination for forming a clear coating film for painted surfaces on metal, glass, ceramic, fiberglass, wood, and waxed substrates. A Treated silica particle layer forms a contact angle of 165 degrees and a rolling angle of 2 degrees. Despite various advantages of the composition, it has the disadvantage that its durability against vibration and bond strength with the surface of the substrate, and is not suited for automotive applications.\nThus, there is a need for a cost effective super hydrophobic and antistatic composition that can be successfully used for all kinds of automobile wheels, overcoming the above mentioned problems of the currently available super hydrophobic coating compositions, which reduces the accumulation of dust, soil and grime on vehicle wheels without the use of harsh chemicals."}
{"text": "There is a number of known illumination devices. Some of the devices are simple incandescent lamps, fluorescent lamps, LED (Light Emitting Diode) based lamps and others. In light emitting diodes (LEDs) based devices, light is produced by one or more LEDs organized to deliver a desired level of illumination on a defined surface area. LEDs are more efficient than incandescent or fluorescent lamps, e.g., LEDs produce up to 100 lumens per watt of consumed electrical energy vs 10-50 lumens per watt produced by incandescent or fluorescent lamps.\nOne of the problems related to LEDs based is reduction in efficiency as operating currents rise, making the devices too hot to power in large-scale applications.\nU.S. Pat. Nos. 5,535,238, 6,791,259, 7,086,765, 7,281,823, 8,742,654, 8,625,097, 8,593,956, 8,948,564 and US Patent Application Publications 2015/0219315, 2016/0061392, disclose different configurations of LEDs based illumination devices.\nHowever, the life constantly requires brighter and more efficient light sources or devices using less energy than the existing light devices. Use of laser diodes as excitation source to excite different inorganic phosphors in illumination devices is believed to be alternative means for high-power white light instead of the traditional LEDs. The more directional nature of the laser diode beams, as compared to LEDs, may be utilized to avoid some optical inefficiencies. The laser diode based lighting devices are higher in efficiency, although in some cases they require excitation by blue laser diodes and a combination of different phosphors.\nPatent Cooperation Treaty WO2005/107240, WO2007/073496 and WO 2007/107420 and U.S. Pat. Nos. 8,625,097, 8,948,564 and US Patent Application Publication No. 20140334126, disclose different configurations of laser diodes based illumination devices."}
{"text": "This invention relates to a method of synthesizing musical tones.\nA musical tone or natural tone produced by an acoustic or natural musical instrument can be expressed by a Fourier series expansion as shown below ##EQU1## where a.sub.n (t) denotes an envelope function of an n-th harmonic component varying with time, .phi.n the phase of the n-th harmonic component, .omega. the angular frequency of a fundamental tone, and N the number of partial frequency components included in the musical tone.\nAs will be evident from equation (1), the amplitudes of spectral components composing a musical tone vary independently with time. That is, a natural tone or actual musical instrument tone is characterized by a dynamic audio spectrum.\nThe simulation of a natural tone with the dynamic audio spectrum can be electronically realized by the use of N envelope function generation circuits, N multipliers and N sinusoidal signal generation circuits. However, since the value of N is relatively large, the actual realization is very difficult.\nIn an existing music synthesizer, a cut-off frequency or frequencies of a voltage-controlled filter is controlled or modulated by a control waveform having a desired shape. The cut-off frequency control of the voltage-controlled filter, however, only vary a frequency range over which higher harmonic components are distributed, and the frequency spectrum of a musical tone being produced substantially remains unchanged. With the music synthesizer it is considerably difficult to obtain dynamic audio spectra."}
{"text": "A conventional ink jet printer has a carriage which can reciprocate in a predetermined direction (primary scanning direction). A recording head is mounted in or on the carriage and has an ink discharge surface facing downward. A recording sheet can be conveyed intermittently at regular intervals in a direction perpendicular to the primary scanning direction. A platen lies below the discharge surface of the recording head and faces the recording head. Ink droplets are discharged from the nozzles of the recording head onto the surface (upper surface) of the recording medium supported on the platen to form an image on the medium.\nThe recording head can discharge ink droplets in a zone including an image recording region. The printer may perform marginless image recording on a recording medium by starting the discharge of ink droplets before the leading end of the medium reaches the edge of the recording region disposed in downstream in the conveying direction, and by performing ink discharge until the trailing end of the medium leaves or passes through the edge of the recording region disposed in upstream in the conveying direction. In this case, there is a problem that the discharged droplets dropping outside both ends of the recording medium adhere to the top surface of the platen which lies below. As a result, these ink droplets stain the under side of the succeeding recording medium supported and sliding on the platen.\nIn order to solve this problem, the printer has been devised as disclosed in U.S. Pat. No. 6,239,817 corresponding to Japanese Patent Application Laid-open No. 2000-118058 (FIGS. 3 and 4A-4C). The recording head of this printer can record an image in an image recording region. The platen of the printer has an upstream wall and a downstream wall, which extend in a primary scanning direction. The upstream and downstream walls stand upright at the edges of the recording region that are upstream and downstream, respectively, in the primary scanning direction. The space between the walls is open upwardly, and an ink absorber is placed in the opened space. The upstream and downstream walls have wavy ribs formed on the top ends thereof, and protrude from the top ends, respectively.\nAccording to the structure of the platen disclosed in U.S. Pat. No. 6,239,817 corresponding to Japanese Patent Application Laid-open No. 2000-118058, however, while a recording medium is passing over the platen, the medium needs to be out of contact with the top side of the ink absorber which is placed in the platen. Therefore, it is necessary to place the ink absorber so that its top side does not protrude over the top ends of the wavy ribs. This makes it necessary to insert, to the platen, an ink absorber which size and shape are strictly controlled. As a result, the incorporation of the ink absorber and other work take time.\nIf the opening space in the platen is large in area, the capacity of ink absorber can be enhanced. Conceivably, it is possible to make this opening space large in area only by increasing the distance between the upstream and downstream walls of the platen in the conveying direction. If the distance between the walls is long, a front end portion of a soft recording medium is likely to hang down when the leading end of the medium has passed over the wavy ribs of the upstream wall and is positioned in the opening space, namely, when the leading end portion is supported in the form of a cantilever by these ribs. Likewise, if this distance is long, a rear end portion of a soft recording medium is likely to hang down when the trailing end of the medium has passed over the wavy ribs of the upstream wall and is positioned in the space, namely, when the rear end portion is supported in the form of a cantilever by the wavy ribs of the downstream wall. In this case, the end portions of the recording medium are likely to come into contact with the top side of the ink absorber in the opening space, so that the under side of the medium is likely to be stained with ink.\nIn order to keep a recording medium out of contact with the ink absorber in the platen space, the vertical distance from the top side of the absorber to the top ends of the wavy ribs may be long. In this case, when the recording head discharges ink droplets erroneously without a recording medium lying on the platen during image recording, the droplets float in the form of ink mist over the absorber. When a recording medium enters the recording region, one or both sides of the medium may be stained by the ink mist."}
{"text": "A microscopic circuit pattern is formed on a surface of a semiconductor wafer, a glass substrate, and the like through a photolithography step or etching step. The chemical liquid supply device has been used in the photolithography step to apply chemical liquid such as photoresist solution to the surface of the wafer or glass substrate, whereby resist solution stored in a chemical liquid tank is sucked up by a pump, passes through a filter and the like, and is applied to a processed object such as the wafer from a nozzle, as described in Japanese Patent Laid-Open Publication No. 9-299862.\nA resist processing device called also a spin coater or resist coater has been used to apply the resist solution to the wafer. The resist processing device includes: a chuck attracting a back surface of the wafer and rotating the wafer; and a nozzle for dropping the resist solution on the wafer. The nozzle reciprocates between an applying position and an evacuating position by an arm, and the resist solution is supplied to the nozzle through a supply tube from the pump provided on a main body side of the resist processing device. An application valve provided on the main body side of the resist processing device is operated when dropping the resist solution from the nozzle."}
{"text": "A device may be provided to sense an environment's present temperature. For example, a temperature probe might be used to determine a current operating temperature inside a boiler, a jet engine, or any other industrial system. One such device, referred to as a “thermistor,” uses a structure whose resistance varies in a pre-determined manner over a range of temperatures. In this case, a voltage across and/or a current flowing through the structure may be measured to determine the temperature of the structure.\nSome thermistors, however, may have impractically small gain characteristics (e.g., the change in resistance may be too small over the range of temperatures that need to be sensed). Moreover, oxidation and other thermal effects may reduce the reliable life of a thermistor when operating at elevated temperatures or limit the operating range of the thermistor. For example, the behavior of the thermistor may begin to drift after prolonged use, especially at relatively high operating temperatures.\nAccordingly, it may be desirable to develop a thermistor that has appropriate gain characteristics over a suitable range of temperatures, including substantially high temperatures. It also may be advantageous to develop a thermistor that is robust and reliable and that can be fabricated at a reasonable cost."}
{"text": "As an optical network system including three or more stations according to wavelength multiplexing in wavelength multiplexing optical transmission, an optical network system using a wavelength branching apparatus (or diverging device) having optical adding/dropping functions (OADM: Optical Add Drop Multiplexer) has been proposed. Due to the OADM, a signal can be arbitrarily inserted (added) or extracted (dropped) in units of wavelength of light.\nIn addition, unlike a fiber branching apparatus, in the wavelength branching apparatus having the OADM, in order to perform the wavelength adding/dropping, a signal band in one optical fiber is divided for communication among three stations. In this case, there is a problem in that, if a cable fault occurs in an arbitrary branch, a quality of a signal communicating between remaining two branches is deteriorated. This problem will be described in detail. In a case where a cable fault occurs, a signal which is to be originally inserted is not multiplexed. Therefore, the number of channels of main signals is decreased. Since a power of a relay is almost constant, a decrease in number of channels of main signals leads to an increase in channel power. In addition, there is a problem in that transmission characteristic according to a nonlinear effect of a transmission line is deteriorated due to the increase in channel power.\nVarious techniques have proposed in order to prevent the deterioration in signal quality according to the cable fault.\nFor example, in a technique disclosed in PTL 1, a dummy light beam for compensating for optical level is disposed in a transmission band (for example, a band L1 in the case of communication between A and B); an interval of cable disconnection is detected from each of terminal station apparatuses A, B, and C in the entire network system; the dummy light beam which is optimized in terms of signal quality is controlled.\nIn a technique disclosed in PTL 2, the deterioration in signal quality is prevented by switching an uplink signal and a downlink signal by using a matrix switch during cable disconnection."}
{"text": "This invention relates to an apparatus for making highly stable and reliable contact, suitable for radio frequency circuit operation, with a component such as a varactor diode or Gunn diode. A prior art assembly which has been used to make contact with such diodes includes a mounting plate assembly and a threaded contact which can be screwed into position to bear against the diode and thus make electrical contact with it. One of the problems with this design is that the threaded contact can stress the diode, breaking it. Another problem is the potential loss of electrical contact with the diode due to thermal expansion and contraction."}
{"text": "Generally described, computing devices utilize a communication network, or a series of communication networks, to exchange data. Companies and organizations operate computer networks that interconnect a number of computing devices to support operations or provide services to third parties. The computing systems can be located in a single geographic location or located in multiple, distinct geographic locations (e.g., interconnected via private or public communication networks). Specifically, data centers or data processing centers, herein generally referred to as a “data center,” may include a number of interconnected computing systems to provide computing resources to users of the data center. The data centers may be private data centers operated on behalf of an organization or public data centers operated on behalf, or for the benefit of, the general public.\nTo facilitate increased utilization of data center resources, virtualization technologies may allow a single physical computing device to host one or more instances of virtual machines that appear and operate as independent computing devices to users of a data center. With virtualization, the single physical computing device can create, maintain, delete, or otherwise manage virtual machines in a dynamic manner. In turn, users can request computer resources from a data center, including single computing devices or a configuration of networked computing devices, and be provided with varying numbers of virtual machine resources.\nIn some scenarios, the physical computing devices or instances of a virtual machine may be configured according to a number of virtual machine instance types to provide specific functionality, often referred to as provisioning. For example, various computing devices may be associated with different combinations of operating systems or operating system configurations, hardware resources (actual or virtualized) and software applications to enable a computing device to provide different desired functionalities, or to provide similar functionalities more efficiently. These customized configurations are often provided within a device image, which a computing device may process in order to implement the desired software configuration and which can be standardized among a set of host computing devices. Once the initial provisioning has occurred, however, at least a portion of the machine image (e.g., the operating environment, software applications, or configurations) may evolve due to being dynamic in nature. Accordingly, the resulting, evolved machine image on a host computing device may diverge from the initial machine image provisioned on the host computing device. For example, some portion of the machine image may correspond to log files and locally persisted state information that would be dynamic and likely different between different host computing devices or virtual machine instances provisioned from the same initial image. As such, it can be difficult for service providers, such as a data center service provider, to utilize the same standardization in providing modifications or upgrades to the machine image data."}
{"text": "This invention relates generally to techniques and apparatus for image sensing, and more particularly to a method and apparatus for obtaining a digital measurement of a photo-current.\nA major problem in imaging is the limited dynamic range of image sensor arrays. As a result, some portions of the image may be saturated while other portions are under exposed. A typical photodiode optical detector, such a charge-coupled device (CCD), consists of an array of pn junction photodiodes. Each photodiode has a capacitance associated with it, and when light falls on the detector the resulting photo-current charges the capacitance. The charge is thus the time integral of the light intensity falling on the detector. The CCD periodically and sequentially switches the charge to a video line, resulting in a series of pulses, which can be converted into a voltage signal representing the light pattern incident on the array. If the integration time is too long, the device will saturate. If the integration time is too short, the voltage signal will be lost in the noise of the device.\nOne approach to adjusting sensitivity is therefore to adjust the integration time according to the intensity of light falling on the array. This can be done globally for the whole array by adjusting the switching rate, for example. This allows the array to adjust for various light levels (as in adjusting the exposure time of an analog camera), but does not increase the dynamic range within a single image.\nAnother approach is to measure the time taken for the capacitor voltage to reach a reference level. An example of such a system is given in xe2x80x9cIntensity Mappings Within The Context Of Near-Sensor Image Processingxe2x80x9d, by Anders xc3x85strxc3x6m, Robert Forchheimer and Per-Erik Danielsson, IEEE Transactions on Image Processing, Vol. 1, No. 12, December 1998. A large time indicates a low intensity of light falling on the array, while a short time indicates a high intensity. The voltage reference level is common to the whole array, but may be adjusted according to the overall intensity falling on the array. This approach is computationally intensive, since the integration time is measured by interrogating the output from a comparator at predetermined intervals. These intervals need to be very close together for accurate measurement. This approach also requires each pixel to have a readout circuit. In addition, each sensor element requires a precharge transistor.\nA further device called a xe2x80x9cLocally Autoadaptive TFA Sensorxe2x80x9d is described in xe2x80x9cHigh Dynamic Range Image Sensors In Thin Film On Asic Technology For Automotive Applicationsxe2x80x9d, by M. Bxc3x6hm et al., in Advanced Microsystems for Automotive Applications, D. E. Ricken, W. Gessner (Eds), Springer Verlag, Berlin, pp 157-172, 1998. In this device, the integration time is adapted for each individual pixel according to the local illumination intensity. The capacitor voltage is compared after an initial integration time to a reference voltage that is slightly less than half of the saturation voltage. If the voltage is less than the reference voltage, the integration time is doubled. This process is repeated until the reference voltage is exceeded. Thus, saturation of the capacitor is avoided and the integration time is extended to avoid under exposure. The final voltage and the integration time are together used to determine the illumination intensity. This approach requires circuitry to compare and adjust the integration time and circuitry to measure the voltage at the end of the integration time.\nIn light of the foregoing discussion, it can be seen that there is still a need in the art to have image sensor arrays resulting in an image representative of the intensity distribution of an object such that portions of the image are not saturated or underexposed."}
{"text": "The use of stress-bearing materials for bone replacement is well known. A large number of designs have been used to replace missing or diseased portions of the bone structure which are stress-bearing, such as bone shafts, joints, and tooth roots. These designs include artificial shafts, joints, and associated devices intended to mimic the functions of the human skeletal system. The designs can be quite elaborate. See, for example, U.S. Pat. No. 3,820,167 which disclosed a design for an artificial hip joint.\nShafts and other prostheses for use as the stress-bearing portion of the bone replacement are typically constructed of metal or metal alloys. Such metal pins or artificial joints are constructed of suitable inert metals such as titanium, stainless steels, or alloys of metals such as cobalt, chromium and molybdenum. The metallic pin or joint is sometimes provided with an oxide coating in order to prevent corrosion and instability.\nIt has been observed in the use of such metallic implants that it is frequently necessary to provide a mechanism to fix its position with respect to adjacent bone and prevent it from shifting in place. Attempts have been made to provide this fixation using a cement. USSR patent application 733,665 published 15 May 1980, discloses an implantfixing cement comprising a mixture of cartilage, collagen (of unspecified origin), and the patient's blood, which is used to secure an implanted pin in place.\nCoating a prosthesis with a material (\"bioglass\" ) which purportedly encourages bone growth, thus securing the implant with the ingrowth, was disclosed in French application Publication No. 2,383,656, published 13 Oct. 1978. It has also been attempted to use wire mesh surrounding the stress-bearing member of the prosthesis to provide sufficient flexibility to obtain a firm fit, in some cases aided by additional metal strips (e.g., EPO Publication 071,242, published 9 Feb. 1983). In connection with this mesh, a biodegradable, biocompatible material (for example, collagen of unspecified origin) is employed, apparently to prevent damage to the wire mesh during insertion of the implant. The resulting implant is eventually secured by additional bone growth. Indeed, research has consistently shown that prostheses which have porous metal or porous ceramic surfaces are better secured by bony growth into the porous surface.\nAn attempt has been made to stimulate this bony ingrowth by providing a bone particle coating. U.S. Pat. No. 3,918,100 describes prostheses made of various metals which are coated using RF sputtering with bone particles in a vacuum. In this disclosure, aluminum oxide substrates are coated with a powder made of lye-treated bone chips. The coated prostheses are said to encourage living bone to grow onto the implant, and the coating finally to be absorbed by the surrounding tissue.\nNone of the foregoing achieves a satisfactory solution. The materials used are often immunogenic and incapable of inducing or stimulating surrounding bone growth at a sufficient rate to prevent damage due to shifting of the prostheses. The present invention provides prostheses which simulate the ingrowth of bone at sufficient rates to obtain a more satisfactory result."}
{"text": "The present invention relates to suspension fuel compositions for use in electrochemical fuel cells, a method of producing electricity with the suspension fuel compositions, and a fuel cell using the suspension fuel compositions to generate electricity.\nA fuel cell is a device that converts the energy of a chemical reaction into electricity. Among the advantages that fuel cells have over other sources of electrical energy are high efficiency and environmental friendliness. Although fuel cells are increasingly gaining acceptance as electrical power sources, there are technical difficulties that prevent the widespread use of fuel cells in many applications, especially mobile and portable applications.\nA fuel cell produces electricity by bringing a fuel into contact with a catalytic anode while bringing an oxidant into contact with a catalytic cathode. When in contact with the anode, the fuel is oxidized at catalytic centers to produce electrons. The electrons travel from the anode to the cathode through an electrical circuit connecting the electrodes. Simultaneously, the oxidant is catalytically reduced at the cathode, consuming the electrons generated at the anode. Mass balance and charge balance are preserved by the corresponding production of ions at either the cathode or the anode and the diffusion of these ions to the other electrode through an electrolyte with which the electrodes are in contact.\nA common type of fuel cell uses hydrogen as a fuel and oxygen as an oxidant. Specifically, hydrogen is oxidized at the anode, releasing protons and electrons as shown in equation 1:\nH2xe2x86x922H++2exe2x88x92xe2x80x83xe2x80x83(1) \nThe protons pass through the electrolyte towards the cathode. The electrons travel from the anode through an electrical load and to the cathode. At the cathode, the oxygen is reduced, combining smith electrons and protons produced from the hydrogen to form water as shown in equation 2:\nO2+4H++4exe2x88x92xe2x86x922H2Oxe2x80x83xe2x80x83(2) \nAlthough fuel cells using hydrogen as a fuel are simple, clean and efficient, the extreme flammability of hydrogen and the bulky high-pressure tanks necessary for storage and transport of hydrogen mean that hydrogen powered fuel cells are inappropriate for many applications.\nIn general, storage, handling and transport of liquids are simpler than for gases. Thus, liquid fuels have been proposed for use in fuel cells. Methods have been developed for converting liquid fuels such as methanol into hydrogen, in situ. These methods are not simple, requiring a fuel pre-processing stage and a complex fuel regulation system.\nFuel cells that directly oxidize liquid fuels are the solution to this problem. Because the fuel is directly fed into the fuel cell, direct liquid-feed fuel cells are comparatively simple. Most commonly, methanol has been used as the fuel in these types of cells, as it is cheap, available from diverse sources and has a high specific energy (5020 Ampere hours per liter).\nIn direct-feed methanol fuel cells, the methanol is catalytically oxidized at the anode, producing electrons, protons and carbon monoxide, as shown in equation 3:\nCH3OHxe2x86x92CO+4H++4exe2x88x92xe2x80x83xe2x80x83(3) \nCarbon monoxide tightly binds to the catalytic sites on the anode. The number of available sites for further oxidation is reduced, reducing power output. One solution to this problem is to use anode catalysts, such as platinum/ruthenium alloys, which are less susceptible to CO adsorption. Another solution is to introduce the fuel into the cell as an xe2x80x9canolytexe2x80x9d, a mixture of methanol Keith an aqueous liquid electrolyte. The methanol reacts with water at the anode to produce carbon dioxide and hydrogen ions, as shown in equation 4:\nCH3OH+H2Oxe2x86x926H++CO2+6exe2x88x92xe2x80x83xe2x80x83(4) \nIn fuel cells that use anolytes, the composition of the anolyte is an important design consideration. The anolte must have both a high electrical conductivity and high ionic mobility at the optimal fuel concentration. Acidic solutions are most commonly used. Unfortunately, acidic anolytes are most efficient at relatively high temperatures, temperatures at which the acidity can passivate or destroy the anode. Anolytes with a pH close to 7 are anode-friendly, but have an electrical conductivity that is too low for efficient electricity generation. Consequently, most prior art direct methanol fuel cells use solid polymer electrolyte (SPE) membranes.\nIn a cell using a SPE membrane, the cathode is exposed to oxygen in the air and is separated from the anode by a proton exchange membrane that acts both as an electrolyte and as a physical barrier preventing leakage from the anode compartment wherein the liquid anolyte is contained. One membrane commonly used as a fuel cell solid electrolyte is a perfluorocarbon material sold by E. I. DuPont de Nemours (Wilmington, Del.) under the trademark xe2x80x9cNafionxe2x80x9d. Fuel cells using SPE membranes have a higher power density and longer operating lifetimes than other anolyte-based fuel cells.\nA practical disadvantage of SPE membrane fuel cells arises from the tendency of high concentrations of methanol to dissolve the membrane and to diffuse through it. As a result, a significant proportion of methanol supplied to the cell is not utilized for generation of electricity, but either is lost through evaporation or is oxidized directly at the cathode, generating heat instead of electricity.\nThe problem of membrane penetration by the fuel is overcome by using anolytes with a low (at most 3%) methanol content. The low methanol content limits the efficiency of the fuel cell when measured in terms of electrical output as a function of volume of fuel consumed and raises issues of fuel transportation, dead weight and waste disposal. Further limiting the use of low methanol content anolyte-based liquid feed fuel cells, especially for mobile and portable applications, is the expense and complexity of necessary peripheral equipment for fuel circulation, replenishment, heating and degassing.\nFinally, despite having a high specific energy, methanol is rather unreactive at room temperature, which limits the specific power output of a methanol fuel cell to about 15 milliwatts per square centimeter.\nOther organic compounds, notably higher alcohols, hydrocarbons and acetates, have been proposed as fuels for fuel cells. See, for example, O. Savadogo and X. Yang, xe2x80x9cThe electrooxidation of some acetals for direct hydrocarbons fuel cell applicationsxe2x80x9d, IIIrd International Symposium on Electrocatalysis, Slovenia, 1999, p. 57, and C. Lamy et al., xe2x80x9cDirect anodic oxidation of methanol, ethanol and higher alcohols and hydrocarbons in PEM fuel cellsxe2x80x9d, IIIrd International Symposium on Electrocatalysis, Slovenia, 1999, p. 95. Most of these candidates have shown very little promise, because of low electrochemical activity, high cost, and, in some cases, toxicity.\nInorganic water-soluble reducing agents, such as metal hydrides, hydrazine and hydrazine derivatives also have been proposed as fuels for fuel cells. See, for example, S. Lel, xe2x80x9cThe characterization of an alkaline fuel cell that uses hydrogen storage alloysxe2x80x9d, Journal of the Electrochemical Society vol. 149 no. 5 pp. A603-A606 (2002), J. O\"\"M. Bockris and S. Srinivasan, Fuel Cells: Their Electrochemistry, McGraw-Hill. New York, 1969, pp. 589-593, and N. V. Korvin, Hydrazine, Khimiya, Moscow, 1980 (in Russian), pp. 205-224. Such compounds have high specific energies and are highly reactive.\nOne such compound is NaBH4. In water NaBH4 dissociates to Na+ and BH4xe2x88x92. In a neutral solution, BH4xe2x88x92 is oxidized at the anode according to equation 5:\nBH4xe2x88x92+2H2Oxe2x86x92BO2xe2x88x92+8H++8exe2x88x92xe2x80x83xe2x80x83(5) \nThe greatest disadvantage of hydrogen-containing inorganic compounds as fuel is their decomposition in acid and neutral solutions. For example. BH4xe2x88x92 decomposes according to equation 6:\nBH4xe2x88x92+2H2Oxe2x86x92BO2xe2x88x92+4H2xe2x88x92xe2x80x83xe2x80x83(6) \nIn a basic solution, BH4xe2x88x92 is oxidized at the anode according to equation 7:\nBH4xe2x88x92+8OHxe2x88x92xe2x86x92BO2xe2x88x92+5H2O+8exe2x88x92xe2x80x83xe2x80x83(7) \nThe corresponding reduction of gaseous oxygen at the cathode proceeds according to equation 8:\nO2+2H2O+4exe2x88x92xe2x86x924OHxe2x88x92xe2x80x83xe2x80x83(8) \nMass balance and charge balance are preserved by diffusion of hydroxyl ions from the cathode to the anode via the electrolyte.\nAlthough BH4xe2x88x92 is stable in basic solutions it decomposes on contact with a catalyst such as is present on the anode of a fuel cell, in accordance with equation 6, even when there is no electrical load on the fuel cell. Although the hydrogen gas produced by this reaction also can be oxidized at the anode, in accordance Faith equation 1, the half-reaction represented by equation 7 is much more efficient, energetically, than the half-reaction represented by the combination of equations 1 and 6. In addition, the catalytic decomposition of BH4xe2x88x92 at the anode tends to shorten the anode\"\"s service life.\nThis problem was addressed in PCT Application No. WO/02/054506, which is incorporated by reference for all purposes as if fully set forth herein, by adding an alcohol such as methanol to the basic NaBH4 solution. In addition to serving as a fuel in its own right, such an alcohol inhibits decomposition of hydride species such as BH4xe2x88x92 at the anode. It is believed that the alcohol inhibits decomposition of hydride species at the anode by at least one of two mechanisms. The first mechanism is that adsorption of alcohol molecules to the anode catalytic sites sterically obstructs access of the hydride species to the catalytic sites. The second mechanism is that alcohol molecules solvate the hydride species.\nIntuitively, it would be expected that capacity (measured in Ampere hours) of a fuel cell that runs on hydride fuel would be a linear function of the hydride concentration. For example, the solubility of NaBH4 in 3M KOH is 1.25 moles per liter, and the solubility of NaBH4 in 3M NaOH is 4 moles per liter, so the capacity of a fuel cell that runs on 3M NaOH saturated with NaBH4 would be expected to be four times that of a fuel cell that runs on 3M KOH saturated with NaBH4. Experimentally, this is not the case.\nFIG. 1 shows, schematically, a fuel cell 10 that consists of an electrolyte chamber 12 that is bounded on either side by a cathode 14 and an anode 16 and that contains an electrolyte. Cathode 14 and anode 16 are shown connected by an electrical load 20 and by an ammeter 22 for measuring the electrical current drawn by electrical load 20. On the other side of anode 16 from electrolyte chamber 12 is a fuel chamber 18 that contains a fuel solution. The oxidant is atmospheric oxygen that reaches cathode 14 on the other side of cathode 14 from electrolyte chamber 12. In the specific fuel cell 10 used in the experiments reported herein, the volume of electrolyte chamber 12 was 2 cm3, the volume of fuel chamber 18 was 15 cm3, and the area of each electrodes 14 and 16 was 4 cm2. Cathode 14 was made by screen-printing 20% platinum on activated carbon on waterproof paper. Anode 16 seas made by screen-printing 20% platinum and 10% ruthenium on activated carbon on hydrophilic carbon paper.\nThe capacity of fuel cell 10 was measured using different concentrations of NaBH4 in a 3.3M aqueous NaOH fuel solution in fuel chamber 18 and using a 6M aqueous KOH electrolyte in electrolyte chamber 12. mF, the effective mass of NaBH4 used, was determined as a function of initial NaBH4 concentration using Faraday\"\"s law:                               m          I                =                  CM          Fn                                    (        9        )            \nwhere C is the measured capacity in Ampere hours, F=26.8 Ampere hours per mole is Faraday\"\"s constant, M=38 g/mole is the molecular weight of NaBH4, and n=8 is the number of electrons released per BH4xe2x88x92 anion in equation 7. The results are plotted in FIG. 2. mF increases with increasing initial NaBH4 concentration, but not linearly. The higher the initial NaBH4 concentration, the less efficiently the NaBH4 is used. Furthermore when the NaBH4 content of the fuel solution exceeded about 50 grams per liter, there was intensive fuel decomposition at anode 16. This, in turn, led to active gas liberation and foam formation, anode process pulsation and gradual destruction of anode 16. Increased initial concentration of NaBH4 also promoted crossover of NaBH4 through anode 16 via the electrolyte to cathode 14.\nThere is thus a widely recognized need for, and it would be highly advantageous to have, a fuel composition for fuel cells that allows a hydride fuel to be used to its full capacity.\nAccording to the present invention there is provided a fuel composition including: (a) a solvent; (b) a first portion of a first fuel, dissolved in the solvent; and (c) a second portion of the first fuel, suspended in the solvent.\nAccording to the present invention there is provided a method of generating electricity including the steps of: (a) providing a fuel cell including a cathode and an anode; (b) contacting an oxidizer with the cathode; and (c) contacting a fuel composition with the anode, the fuel composition including: (i) a solvent, (ii) a first portion of a fuel, dissolved in the solvent, and (iii) a second portion of the fuel, suspended in the solvent.\nThe present invention is a fuel composition for fuel cells in which a first fuel is stored in two forms. A first portion of the first fuel is stored in solution in a solvent. A second portion of the first fuel is stored in suspension in the solvent. The effective concentration of the first fuel is the concentration of the first fuel in solution, and this concentration is kept low enough to preclude undesirable side effects such as decomposition of the first fuel at the anode and destruction of the anode. As the dissolved first fuel is used up, it is replaced by dissolution of the suspended first fuel. The effective mass of the first fuel is close to the total mass of the two portions of the first fuel.\nPreferably, the solvent is a polar solvent such as water. Preferably, the concentration of dissolved first fuel is the saturated concentration of the first fuel in the solvent. During the course of the operation of a fuel cell, as the dissolved first fuel is consumed, the suspended first fuel replaces the dissolved first fuel in solution and so maintains the dissolved portion of the first fuel at its saturated concentration.\nPreferably, the first fuel is a salt whose anion is a product of a reduction half-reaction, in the solvent, that has a standard reduction potential more negative than the standard reduction potential of a hydrogen electrode in the solvent. For example. BH4xe2x88x92, the anion of NaBH4, is the anion produced in the reduction half-reaction (in water)\nHBO3xe2x88x92+5H2O+8exe2x88x92xe2x86x928OHxe2x88x92xe2x80x83xe2x80x83(10) \nwhich has a standard reduction potential of xe2x88x921.24 volts.\nPreferably the first fuel is a hydride such as LiAlH4, NaBH4, LiBH4, (CH3)3NHBH3, NaAlH4, NaCNBH3, CaH2. LiH, NaH or KH. Most preferably, the first fuel is NaBH4. Other preferred first fuels include Na2S2O3, Na2HPO3. Na2HPO2. K2S2O3, K2HPO3, K2HPO2, NaCOOH and KCOOH, which, like the hydrides, are salts whose anions have standard reduction potentials in water that are more negative than the standard reduction potential of a hydrogen electrode in water. For solvents generally, the preferred fuels for any specific solvent include salts whose anions have standard reduction potentials in that solvent that are more negative than the standard reduction potential of a hydrogen electrode in that solvent. Preferably, the first fuel constitutes between about 0.1% and about 80% of the fuel composition by weight. Most preferably, the first fuel constitutes between about 5% and about 25% of the fuel composition by sleight.\nOptionally, the fuel composition of the present invention also includes an alcohol, for example methanol, ethanol, propanol, butanol, pentanol, hexanol, ethylene glycol or glycerol. Preferably, the alcohol constitutes between about 0.1% and about 50% of the fuel composition by weight. Most preferably, the alcohol constitutes between about 1% and about 25% of the fuel composition by weight. The alcohol serves four functions:\n1. The alcohol is a second fuel that is oxidized along with the first fuel at the anode of the fuel cell.\n2. The alcohol controls the solubility of the first fuel in the solvent, to ensure that the saturated concentration of the first fuel is not too high.\n3. As in WO/02/054506 with respect to NaBH4, the alcohol inhibits the decomposition of the first fuel at the anode of the fuel cell.\n4. The alcohol stabilizes the suspension by being present, in the solution of the first fuel in the solvent, in a proportion that makes the density of the solution substantially equal to the density of the suspended portion of the first fuel, so that the suspended portion of the first fuel neither precipitates nor floats, but remains suspended.\nThe scope of the present invention also includes the use of any suitable additive for any of these four purposes, but alcohols are the preferred additives.\nPreferably, the fuel composition of the present invention includes an additive for stabilizing the dissolved portion of the first fuel in the solvent. Preferably, this additive is an alkali such as LiOH, NaOH or KOH, or a basic salt. Preferably, this additive is present in the solvent in a concentration between about 0.1 mole/liter and about 12 mole/liter. Most preferably, this additive is present in the solvent in a concentration between about 0.2 mole/liter and about 5 mole/liter.\nOsborg, in U.S. Pat. No. 4,081,252, teaches a fuel composition, for combustion rather than for use in a fuel cell, that, similar to the present invention, includes a xe2x80x9chydrogen carrierxe2x80x9d such as hydrazine, a hydrazine derivative or an inorganic borohydride, that, according to the abstract of the patent, may be dissolved or suspended in a base fuel. All the examples presented by Osborg, however, are of hydrogen carriers that are dissolved in the base fuel. There is no indication in Osborg of any utility to both dissolving and suspending a hydrogen carrier in the base fuel.\nThe scope of the present invention also includes a fuel cell that is fueled by the fuel composition of the present invention, as well as a method of generating electricity using such a fuel cell."}
{"text": "In the front end manufacture of CMOS semi-conductor devices, a passivation film such as SiN is formed on the gate electrode of each MOS transistor. This SiN film deposited on top and on side surface of the gate electrodes (such as polycrystalline silicon or metallic layers) in order to increase the breakdown voltage of each transistor.\nAttempts have been made to reduce the temperature deposition of such SiN, to reach a temperature which is not higher than 400° C.\nHowever, SiN films deposited at temperatures below 400° C. have usually poorer film qualities. In order to overcome this issue, it has been proposed to use SiO2 films to reinforce SiN film properties (“dual spacer”) and thereby make effective electrical barrier layers to improve significantly the device performances. Also, SiO2 films are used as STI (shallow trench insulation), inter layer dielectric (ILD) layers, passivation layers, etch-stop layers and attempts are made to find a deposition process of these SiO2 layers at low temperature, i.e. below 400° C. In the specific case of dual spacer applications, the deposition of a very (20-50 A thick) thin films performed at low deposition temperatures (300° C. at the most), should not lead to the oxidation of the metal electrode and should be perfectly uniform all along the gate, and an atomic layer deposition process is the most suitable such a requirement. As far as the STI applications are concerned, conformal films should be deposited with high deposition rate (several hundred A per minute) below 500° C.\nDeposition of silicon oxide films made from silane and oxygen at low temperature using a PECVD reactor have failed up to now for these applications, due to the incorporation of Si—H bonds into the SiOx film thereby obtained, which may originate with the presence of hydrogen in the silane gas introduced as a precursor gas. This Si—H bonds then probably readily react with an oxygen source in the environment to produce the Si—Ohio bond.\nThe presence of such Si—OH bonds increases the risk of having transistors with an increased leakage current, and therefore a reduced breakdown voltage of the related transistors.\nThe inventors believe that the presence of a large number of hydrogen atoms bonded to the silicon atom in the Si precursor and of oxygen or oxygen containing gas to generate the SiO2 film, probably also generates moisture (H2O) formation which then reacts with Si to make SiOH.\nThe inventors also believe that the silicon containing compound shall have preferably a high purity regarding hydrogen containing residues to limit H2O generation, preferably it shall contain less than 100 ppm of H2 or H containing compounds."}
{"text": "1. Field of the Invention\nThe present invention relates to a driving force transmission apparatus having a toothed clutch that transmits a driving force received from a driving source, and a sheet conveyance apparatus.\n2. Description of the Related Art\nAn image forming apparatus, such as a copier, a printer, a facsimile machine or a multifunctional apparatus having all these functions, mounts a thermal fixing device and a sheet conveyance mechanism, such as a recording material conveying belt mechanism. This sheet conveyance mechanism includes a sheet conveyance roller (or belt), which forms a sheet conveyance nip portion; a drive motor, for driving a sheet conveyance roller; and a driving force transmission mechanism, which transfers the driving force produced by a driving motor to the sheet conveyance roller.\nWhen the image forming apparatus is abnormally halted because of a sheet jam in the sheet conveying mechanism and the sheet is caught at the conveyance nip portion, the sheet must be manually removed by a user, i.e., the jam must be cleared by the user. While the jam is being cleared, the sheet conveyance roller is rotated by friction between it and the sheet. At this time, because of the operating arrangement of the image forming apparatus, a driving force generated by the rotation of the sheet conveyance roller is transmitted to the driving motor via the driving force transmission mechanism, that is, a load produced by the rotation is imposed on the gear of the driving force transmission mechanism and thence on the driving motor. Therefore, removing the sheet is difficult and the sheet may sometimes be torn.\nIn order to reduce such a load, a driving force transmission mechanism having a function in which it is able to stop transmitting a driving force may be provided between the driving motor and the sheet conveyance roller. A toothed clutch is one example of such a driving force transmission mechanism. A toothed clutch includes a driving-side ratchet and a driven-side ratchet that are biased in the direction in which they are coupled. Thus, when the driving motor is rotated in the sheet conveying direction, the driving-side ratchet and the driven-side ratchet engage each other to transmit the driving force produced by the driving motor to the sheet conveyance roller. On the other hand, when a jam is being cleared and the sheet conveyance roller is rotated in the direction opposite to that in which a sheet is conveyed, the pawl of the driven-side ratchet slides over the pawl of the driving-side ratchet, so that the driven-side ratchet is not rotated and transmission of the driving force from the sheet conveyance roller to the driving motor is prevented. In this case, since the load imposed on a jammed sheet is comparatively small, a user can easily remove the sheet from the sheet nip portion.\nHowever, for the conventional art described above, there are the following problems.\nA user must rotate the driven-side ratchet forward a number of times when clearing a jam where a sheet is caught at the conveyance nip portion of the sheet conveyance roller. Therefore, the pawl of the driven-side ratchet must slide over the pawl of the driving-side ratchet many times. And since the driven-side ratchet and the driving-side ratchet are biased forward by a spring, in the direction in which they are coupled, the ratchet pawls forcefully strike each other each time the pawl of the driven-side ratchet slides over the pawl of the driving-side ratchet, and an impulsive sound (an action sound) is produced. It is assumed that this impulsive sound will frequently cause user discomfort.\nAs means for preventing the generation of such an impulsive sound, there is an arrangement whereby, when a user opens the cover of an apparatus to clear a jam, a lever that interlocks with the cover reversely impels the pawl of a driven-side ratchet separating, against the biasing force exerted by an biasing spring, away from the pawl of a driving-side ratchet (see Japanese Patent Application Laid-Open No. H10-142996). According to this arrangement, since the ratchet pawls do not forcefully strike each other while the user is removing the jam, the occurrence of the impulsive sound can be prevented. However, since a lever and an interlocking mechanism for securing the main body of the apparatus to the lever must be additionally provided, the number of parts is increased.\nA pressure release mechanism is also proposed whereby, when a sheet jam has occurred, a drive motor is rotated in reverse, in the direction opposite that in which a sheet is normally conveyed, to remove the pressure imposed on a conveyance nip portion (see Japanese Patent Application Laid-Open No. 2000-029347). According to this arrangement, since no pressure is imposed on the conveyance nip portion while the jam is being cleared, the process can be easily performed.\nHowever, the pressure release mechanism, for example, may not work, as in a case wherein power to the image forming apparatus is turned off during the image forming processing. At this time, while a user is clearing the jam, the ratchet pawls would strike each other and produce the above described impulse sound.\nFurther, according to the arrangement whereby the pressure imposed on the conveyance nip portion is released by rotating the driving motor in reverse, in the direction opposite to that in which a sheet is conveyed, the driving-side ratchet pawl will slide over the driven-side ratchet pawl each time the pressure release procedure is performed. And thus, an impulse sound will still be generated by the ratchet pawls."}
{"text": "1. Field of the Invention\nThe present invention relates to personal communications systems, including, but not limited to cellular telephones, personal digital assistants, pagers, and internet appliances.\n2. Problems in the Art\nThe present invention relates to personal electronic devices and personal communication systems and many problems associated with such devices. One of the most recognized problems relates to hands-free communication. In voice communication systems, such as cellular telephones, it is more convenient and safer for voice communications to occur without requiring the use of hands. Many attempts have been made at solving this problem, resulting in cords running between headsets and cellular telephones, heavy and bulky headset units, and related problematic solutions. Another problem concerning personal electronic devices and personal communication systems is the increasing complexity of such devices. As technology improves, additional functionality has been given to cellular telephones, personal digital assistants, pagers, internet appliances, and other personal electronic devices. Increasing the complexity of the device generally involves increasing the number of buttons or controls associated with the device or otherwise making a device more hands intensive in operation.\nThus there is an even more pronounced problem when additional functionality is given to a personal digital assistant, cellular telephone, or related personal electronic device when that personal electronic device has been designed to provide additional functionality. The number of controls is increasingly providing greater difficulties in operating the device and there are greater difficulties in accessing the controls or buttons on the device. This is particularly true where some of the controls, such as a touch screen are protected.\nTherefore, it is a primary objective, feature, or advantage of the present invention to provide an apparatus, and system which improve upon the state of the art.\nIt is another objective, feature, or advantage of the present invention to provide a system and apparatus of personal communication.\nIt is another objective, feature, or advantage of the present invention to provide a system and apparatus of hands-free communication.\nIt is a further objective, feature, or advantage of the present invention to provide a system and apparatus having accessible access to user controls.\nIt is a further objective, feature or advantage of the present invention to provide a system and apparatus that provides simplified access to a touch screen display.\nIt is a further objective, feature, or advantage of the present invention to provide a personal electronic device that includes a cellular transceiver and a personal digital assistant.\nIt is another objective, feature, or advantage of the present invention to provide a system and apparatus that includes a cellular transceiver, a personal digital assistant, and a pager.\nIt is a further objective, feature, or advantage of the present invention to provide a personal electronic device with a telephone keypad and a touch screen display, the telephone keypad and the touch screen display being simultaneously accessible.\nThese and other objectives, features, or advantages of the present invention will become apparent from the specification and claims."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, in particular, to a semiconductor device having high electrostatic discharge (ESD) immunity.\n2. Description of the Related Art\nAn ESD element is indispensable in maintaining reliability of an integrated circuit (IC), though it is not directly responsible for a function of the IC. The ESD element refers to an electrostatic discharge element and serves to release static electricity in order to avoid destruction of the IC by the static electricity.\nAccordingly, it is an essential condition that the ESD element itself is not thermally destroyed by the static electricity and is able to protect an internal circuit by dissipating electric charges quickly before the static electricity enters the internal circuit. In order to satisfy this condition, the ESD element is required to have properties of high driving ability and suppressing local heat generation. A generally employed method is to increase a sectional area of the ESD element at a point through which current flows, resulting in inevitable increase in the size of the ESD element. It is thus important to obtain the above-mentioned properties within an ESD element of a small size.\nFurther, in order not to disrupt the normal function of the IC, the ESD element is required to have a breakdown voltage that is not smaller than the absolute maximum rating of the IC. Especially in a case of a high withstanding voltage IC, the above-mentioned local heat generation problem becomes more serious since the ESD protection element must dissipate electrostatic charges while being applied with a voltage that is not smaller than the absolute maximum rating.\nJP 2004-335634 A discloses a structure for preventing melting and destruction of an electrode due to local heat generation of an ESD element. FIG. 2 shows a sectional view of a conceptually illustrated structure. When positive static electricity is injected, for example, to an input pad (PAD), the static electricity flows through a collector electrode 7 to an n+ collector layer 2 and into an N-well electric field relaxation layer 23 for improving withstanding voltage. However, with the N-well electric field relaxation layer 23 having high resistance, electric charges tend to accumulate in the n+ collector layer 2, and the electric field becomes strong to locally generate heat around a boundary (“heat generation area” encircled in FIG. 2) between the n+ collector layer 2 and the N-well electric field relaxation layer 23 on a path having a small distance from the n+ collector layer 2 to an n+ emitter layer 6. The generated heat is conducted to a collector contact region 1 to melt the collector electrode 7. In order to solve this problem, the prior are attempts to prevent the melting and destruction of the collector electrode 7 by extending the distance b from the above-mentioned heat generation area to the collector contact region 1.\nSince the method disclosed in JP 2004-335634 A is not, however, for suppressing local heat generation itself, melting and destruction of the silicon where the heat is generated may occur. Further, extension of the distance b not only increases the size, but also increases resistance of the n+ collector layer 2 and reduces the driving ability of the ESD element itself to increase a possible risk of exposing the internal circuit to a static electricity. A long base length of the ESD element is needed to increase the driving ability, further giving an increase in the size.\nAs described above, in order to obtain sufficient properties without increasing the size of the ESD element, it is important to suppress the local heat generation itself in the ESD element."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an apparatus for mounting items onto vehicles; more specifically, the invention relates to a rigid and adjustable apparatus for supporting gauges or other items from available areas of vehicles, whether open surfaces or tubular members such as handlebars, roll bars, or A-pillars.\n2. Background Information\nSome forms of personal transportation lack desired controls, indications, or accessories. Vehicles such as bicycles or snowmobiles offer little space typically considered suitable for instrumentation. Some manufacturers forego such equipment to keep their costs low. These problems provide opportunities for creative after-market solutions. Because of the expense of design and production, after-market manufacturers try to design products applicable to a wide variety of makes or styles of vehicles. Therefore, it is important that after-market solutions be adjustable, capable of suiting most drivers regardless of the vehicle. If a gauge mount is not adjustable, then it may not permit ergonomically correct positioning. Ideally, the flexibility of orientation and location would help a variety of drivers of a single vehicle to achieve spatial compatibility with the installed indications. Such products should also be effective, inexpensive, and easy for the driver to install. For performance racing vehicles, it is also important for the gauge mount to fixedly hold the gauge in place through high speeds and hard turns.\nMost vehicles feature some exposed, available surface. For example, motorcycles and bicycles commonly have tubular handlebars. Vehicles used in motor sports are likely to provide roll over protection through tubular roll bars or cages; additionally, in some of these vehicles the A-pillar is also an exposed tube. Larger vehicles, such as boats and automobiles may have free surfaces on dashboards or consoles. If such spots are adjacent to the vehicle driver, they can provide an opportunity for mounting a gauge, control, or accessory.\nOne approach, seen mostly on two wheeled vehicles, is to create a support bracket with holes that allows the bracket to be held to the vehicle by existing vehicle bolts that are passed through the bracket. Examples of this technique are disclosed in U.S. Pat. No. 5,855,388 to Brewer and U.S. Pat. No. D 367,418 to Lathrop. These inventions often rely on the common availability of bolts near the driver's view in the center of bicycle or motorcycle handlebars; the bolts may fasten the head of the handlebar assembly, tie into the front forks, or anchor the top of the suspension. This approach is limited in several ways. First, the positioning is limited to the location and availability of the mounting bolt or bolts, possibly inhibiting good spatial compatibility. Second, the installation requires removing a bolt, possibly from the handlebar or front suspension; novice mechanics may be uncomfortable undertaking this task. Third, it is possible that the added thickness of the bracket could impede the proper functioning of the mounting bolt if its threads are no longer able to engage.\nAnother static or fixed approach, common for larger vehicles such as automobiles and boats, is to fix a mounting structure or frame to an available surface; this structure ideally holds and orients the gauge or accessory towards the driver. Some static, after-market gauge mounts are available for placing gauges on the driver side A-pillar, in pods on the surface of dashboards, or adjacent to existing instrument clusters. These solutions often involve a tradeoff between spatial compatibility and vehicle applicability. If the manufacture seeks to optimize ergonomics, then the product is likely to be suitable only for a single or small number of vehicle models. For example, a pod designed to be in the driver's view when placed on the curved dash of a subcompact car is unlikely to be suitable for the straight dash of a large sedan. Some static mounts offer adjustment of the gauge in a single plane, such as pitch. In general, however, static mounts are limited in their ability to achieve spatial compatibility for a wide variety of drivers.\nAlternatively, some devices take advantage of the space available on exposed tubular members, creating brackets that clamp onto features such as handlebars. U.S. Pat. No. 3,843,083 to Angibaud, U.S. Pat. No. D 257,122 to Price, U.S. Pat. No. 5,860,728 to Maglica, and U.S. Pat. No. 6,331,089 to Iteya all seek to mount gauges or accessories to exposed tubes. The problem with this approach is gripping the tube in a rigid manner, while offering flexibility in adjustment.\nA number of these designs use a form of a single split clamp and gauge support. The Angibaud apparatus forms a sort of single split clamp when the single screw is tightened. Similarly, '122 to Price, '089 to Iteya, and several of the embodiments in '728 to Maglica use single split clamps. In some cases, such as in the Angibaud and Price designs, the clamp uses a single fastener to provide the clamping force for both the handlebar and the gauge support; when the screw is loosened to adjust one clamp, the apparatus loses its grip on both single split clamps. Depending on the orientation and shape of the bar, for example, single split clamps can also suffer from uneven distribution of compressive forces, reducing the friction available to hold the gauge in place.\nSome devices incorporate double split clamps. For example, U.S. Pat. No. 5,413,007 to Vernon depicts a double split handlebar clamp with a plurality of machine screws. The Vernon design features enhanced stability from a double split clamp with a fixed mount. However, the gauge may be adjusted only through the rotation of a threaded member. This design limits the range of positions available for the gauge.\nThe design in '728 to Maglica uses a double split arrangement to grip handlebars, and a form of ball joint. This invention centers on a spherical ring “gripping” member, which is a ball or sphere with a cylindrical hollow into which can be inserted a removable flashlight. The mounting of this gripping member offers some ability to rotate for adjustment, but the housing must retain the gripping member within its mount, which limits its rotation when it is “gripping” or holding a cylindrical object like a flashlight. For the intended application of the '728 invention, this is not a serious problem because the angles needed for a flashlight to aid a bicycle driver at night are relatively limited. However, it would pose several problems if one tried to adapt the Maglica arrangement for use with a gauge.\nDifferent drivers, and even different positions of a single driver can require a broad range of adjustment to preserve spatial compatibility of an important gauge or indication. Additionally, items such as gauges typically have a greater diameter to length ratio, which means that the spherical ring “gripping” member and its housing would have to be much larger to accommodate the larger diameter. This increased size could inhibit adjustment, add weight, obscure the view of the gauge, and create wind resistance. For the reasons discussed above, these adaptations are likely to reduce further its functionality.\nIn several embodiments, the Maglica ball is positioned in a double split housing with the additional function of clamping to the handlebar. This approach shares the drawback with single fasteners described above, in that any loosening of the handlebar clamp would similarly loosen the force on the floating ball until it is ultimately released from the housing. Likewise, loosening the housing to adjust the “gripping” member could result in a loss of grip on the handlebars. In another Maglica embodiment, the “gripping” member is held in place by compression of a washer by a threaded ring. In that case, the clamping of the ball is achieved through the friction of an appropriately sized rubber washer.\nNotably, the patents to Iteya, Angibaud, and Price also use clamps and types of ball joints to make their gauge mounts adjustable. The Iteya design, like the Maglica design, is limited in that it involves inserting an item within the spherical ring, restricting its range of motion. Angibaud and Price rely on single split clamps to grip both handlebars and a spherical member; the spherical member is then connected by a neck to the gauge, accessory, or additional mount. Aside from the single fastener weakness described above, these designs need a compressive housing to keep the sphere steady under the torque of a gauge positioned at the end of this neck. The compressive force and friction of the housing must be able to overcome bumps and forces of everyday driving. At the same time, this housing limits the range of motion for gauge adjustment.\nNone of these designs offers both rigidity and a wide range of flexibility. The present invention is directed to an improved gauge mount for overcoming this challenge."}
{"text": "1. Field of the Invention\nThis invention relates generally to fluid injectors. More particularly, this invention relates to a proportional needle control fuel injector having a piezoelectric actuator and an actuator motion amplifying lever.\n2. Description of Related Art\nAccurate control of fuel injection rate shape requires accurate control of the injector needle valve motion. Highly accurate electronic control of an injector needle valve motion may be provided by using solid state actuators such as piezoelectric, electro-strictive, or magneto-strictive actuators to control the needle valve position. To provide a desired valve displacement, conventional solid state actuators require high values of applied current and/or voltage.\nIn an attempt to reduce the magnitude of the applied control current and/or voltage to the solid state actuators, and dimensions of the actuator, some fuel injectors include mechanical or hydraulic motion amplifiers. The use of a hydraulic amplifier requires that some fuel must be wasted to control the amplification. Hydraulic amplifier\"\"s also allow undesirable oscillations of the controlled motion parts\nA piezoelectric actuator stack is advantageous in that a piezoelectric actuator reacts approximately five to ten times faster than a solenoid controlled injector. Additionally, the amount of extension of the piezoelectric actuator is directly proportional to the amplitude of the control current/voltage, applied to the piezoelectric actuator. A proposed actuator with a lever system is also beneficial in that the injected fuel may be accurately controlled without the use of a spill valve circuit. Conventional spill valves maintain a desired pressure within a fuel injector by bleeding off fuel from the fuel passageway when the pressure within the fuel passageway exceeds a predetermined amount. The fuel that is spilled either returns to the fuel supply or is wasted. A spill valve circuit needle control system is inherently inefficient in that the energy that is used to pressurize the fuel is wasted when it is spilled from the fuel system. In addition to the energy expended in increasing the pressure of fuel, energy may also be expended in heating the fuel and this heat may also be lost in a spill valve circuit.\nSolid state actuators provide such an accurate degree of needle positioning that variable geometry atomizing orifices may be effectively used. Variable geometry atomizing orifices enable high quality atomization for all operation conditions of the engine, and accurate control over the amount and rate of fuel being injected.\nIt is an object of the present invention to overcome the disadvantages of the prior art and to provide a proportional needle control injector that compensates for changes in fuel pressure.\nIt is another object of the present invention to provide a proportional needle control injector that compensates for changes in friction forces.\nIt is another object of the present invention to provide a proportional needle control injector that compensates for wearing.\nIt is yet another object of the present invention to provide a proportional needle control injector that compensates for manufacturing tolerances.\nIt is a further object of the present invention to provide a proportional needle control injector that improves injection rate shape control.\nIt is a still further object of the present invention to provide a proportional needle control injector with reduced requirements to control current and/or voltage.\nIt is still another object of the present invention to provide a proportional needle control injector that does not require a control flow.\nIt is yet another object of the present invention to provide a proportional needle control injector that saves deficit radial space.\nIt is another object of the present invention to provide a proportional needle control injector that is hydraulically compensated such that the amplitude of the control current/voltage to the actuator may be reduced and the strength of a spring biasing the needle valve closed may also be reduced.\nThese and other objects of the present invention are achieved by providing an injector body including an injection fuel inlet port, an injection fuel outlet port and a control fuel inlet port. The injector also includes a solid state actuator housed within the injector body, an actuator amplifying lever housed within the injector body and a needle valve housed within the injector body. The needle valve being moveable between a first position closing the injection fuel outlet port and a second position to open the injection fuel outlet port. The actuator amplifying lever is responsive to a dimensional change of the solid state actuator to permit movement of the needle valve from the first position to the second position. The distance between the first and second positions (needle valve travel/lift) is proportional to the applied voltage/current.\nThese and other objects of the present invention may also be achieved by providing a fuel injector having an injector body that includes an injection fuel inlet port and an injection fuel outlet port. The injector also includes within an injector body: a solid state actuator, an actuator amplifying lever, a control piston and a needle valve. The control piston is housed within a control fluid chamber within the injector body which is in fuel communication with the injection fuel inlet port and the injection fuel outlet port. The needle valve is moveable between a first position closing the injection fuel outlet port and a second position to open to injection fuel outlet port. The actuator amplifying lever is responsive to a dimensional change of the solid state actuator to permit movement of the needle valve from the first position to the second position.\nThese and other objects of the present invention may also be achieved by providing a fuel injector having an injector body that includes an injection fluid inlet port and an injection fluid outlet port. The injector body houses a solid state actuator, a needle valve, a biasing means and an actuator amplifying lever. The needle valve is moveable between a first position closing the injection fluid outlet port and a second position to open the injection fluid outlet port. The biasing means biases the needle valve toward the first position. The actuator amplifying lever is positioned between the biasing means and the needle valve and is responsive to a dimensional change of the solid state actuator to permit movement of the needle valve from the first position to the second position.\nOne exemplary embodiment of the present invention provides hydraulic compensation for clearances between the solid state actuator and the needle valve. This embodiment uses a fuel that compensates for wear and manufacturing tolerances.\nIn another exemplary embodiment of the present invention the amplitude of the required current and/or voltage for the solid state actuator may be reduced by using a separate control channel to control forces acting at the top of a control piston linked to a needle valve.\nYet another exemplary embodiment of the present invention may use a pusher pin to contact an amplifying lever and to minimize the size of the fuel injector.\nAn additional exemplary embodiment may include a mechanical amplifying lever for the solid state actuator that engages an integral needle valve and control piston."}
{"text": "Relays may be used in LTE-A to enhance coverage and capacity and offer more flexible deployment options. A relay node (RN) may create new cells, distinguishable and separate from the cells of a donor-evolved NodeB (eNB), where the donor-eNB may support a Un interface to support eNB to RN communication. To any legacy wireless transmit/receive unit (WTRU), an RN may appear as an eNB. That is, the presence of an RN in its communication path to the donor-eNB may be transparent to the WTRU. An RN may be an eNB that has a wireless in-band backhaul link back to the donor-eNB by using an LTE or LTE-A air interface within the International Mobile Telecommunications (IMT) spectrum allocation.\nA two-way communications system with two transceivers and one relay node may take four time slots to complete a message exchange, assuming time division duplex (TDD) half duplex mode and no direct link between two transceivers. It would be desirable to have a method and apparatus for implementing a relay scheme that may use a lower number of time slots and may also have a low level of complexity."}
{"text": "1. Field of the Invention\nThe present invention relates to a composition for improving a taste of a high-intensity sweetener, a method of improving a taste of an edible product containing a high-intensity sweetener, and an edible product in which a taste of a high-intensity sweetener is improved. More specifically, the present invention relates to a taste improving composition which can improve a taste of a high-intensity sweetener by suppressing the exhibition of a strange taste (bitterness) of the high-intensity sweetener, a method of improving a taste of an edible product, and an edible product, particularly a low-calorie edible product, in which a taste of a high-intensity sweetener is improved.\n2. Description of the Related Art\nSugar and an isomerized corn syrup (high fructose corn syrup; HFCS) are sweeteners which have a favorable quality of sweetness and are popular among many people. However, recently, it has been pointed out that the excessive intake of such sugar or an isomerized corn syrup causes obesity or lifestyle diseases. In such circumstances, from health conscious point of view, an attempt has been continually made to obtain a low-calorie food or drink by substituting sugar alcohol or high-intensity sweetener having a sweetness several hundreds times higher than that of sugar for part or most of the sugar or isomerized corn syrup when producing a food or drink.\nHowever, although the high-intensity sweetener has been frequently used as a low-calorie sweetener because the high-intensity sweetener has a sweetness several times to several thousands times higher than that of sugar (sucrose) and can exhibit a desired sweetness in a small addition amount, not only does the sweetness persist so as to deteriorate the sharpness of aftertaste, but also the aftertaste is bad such that a distinctive strange taste or bitterness is exhibited as the aftertaste. In addition, the high-intensity sweetener has a disadvantage, that is a lack of body (robustness). Therefore, the high-intensity sweetener is rarely used alone and several types of high-intensity sweeteners are used generally in combination. For example, a high-intensity sweetener such as aspartame, acesulfame potassium, or sucralose (registered trademark) has been widely used in a diet food or drink because of its high sweetness potency. However, consumers have been familiar with the taste of a so-called regular product obtained by using sugar or an isomerized corn syrup (a fructose-glucose syrup or a glucose-fructose syrup), and therefore, there are many consumers who feel unpleasant to the sweetness of a food obtained by using such a high-intensity sweetener. Many studies have been made for these three types of high-intensity sweeteners, for example, a combination use thereof with trehalose, erythritol, or the like (JP-A-2002-51723), a combination use thereof with dietary fiber (JP-A-2004-41118), and a combination use thereof with an α-glucosylated stevia extract (JP-A-2002-34501) have been disclosed, however, the current situation is that the satisfaction of the preference of consumers who have been familiar with the taste of sugar or a high-fructose corn syrup has not been obtained yet.\nOn the other hand, D-psicose is a rare sugar contained in molasses or the like in a very small amount. Although the sweetness of D-psicose is 60 to 70% of that of sugar (sucrose), D-psicose has almost zero calorie, and also has a high solubility, and therefore, D-psicose is expected to be applied to various foods and drinks. However, in the case of using D-psicose for imparting sweetness to a food or drink, when D-psicose is added until necessary sweetness is imparted, the added amount thereof would be much increased, and gives a much too rich taste or its sweet taste is not felt instantaneously. In consideration of these facts, single use of D-psicose is not realistic, and therefore, a method of supplementing the disadvantage of D-psicose by using a sugar alcohol in combination has been disclosed (JP-T-2008-059623)."}
{"text": "1. Technical Field\nThis disclosure relates to a semiconductor memory device, and more particularly, to a multi path accessible semiconductor memory device for use in a portable communication system.\n2. Description of the Related Art\nIn general, a semiconductor memory device having multiple access ports is called a multi port memory, and in particular, a memory device having two access ports is called a dual port memory. A typical dual port memory is an image processing video memory having a RAM (Random Access Memory) port accessible by a random sequence and an SAM port accessible only by a serial sequence.\nToday, electronic systems are becoming ubiquitous in human life. For example, handheld phone or PDA (Personal Digital Assistant), etc., manufacturers have been produced multi processor systems having multiple processors as shown in FIG. 1 so as to obtain high speed functionality and smooth system operation.\nReferring to FIG. 1, a first processor 10 is connected to a second processor 26 through a connection line L10. A NOR memory 14 and a DRAM 40 are connected to the first processor 10 through buses B1-B3. A DRAM 18 and a NAND memory 20 are connected to the second processor 26 through determined buses B4-B6. The first processor 10 may have a base band processing function or may perform a modulation and/or demodulation of communication signal. The second processor 26 may have an application function of processing communication data or performing play activity known as a game, etc. The NOR memory 14 has a NOR structure cell array configuration and the NAND memory 20 has a NAND structure cell array configuration. Both are nonvolatile memories having transistor memory cells that have floating gates, and are equipped to preserve and store data that should not be lost if power is removed. For example, such data includes device specific code of handheld instruments or configuration data. The DRAMs 40 and 18 function as main memories of processing data of processors 10 and 26.\nHowever, in a multi processor system as shown in FIG. 1, DRAMs are each allocated to each processor. In addition, an interface, such as a UART (Universal Asynchronous Receiver/Transmitter), SPI (Serial Peripheral Interface), and SRAM (Static Random Access Memory), having a relatively low speed are used to access the DRAMs. Thus, it is difficult to obtain a satisfactory data transmission speed and complexity due to size is increased. Furthermore, cost increases because of configuring the memories. In addition, data needed by both processors 10 and 26 is transmitted through a data transmission channel L12, adding a delay of data transmission.\nA scheme to reduce an occupied size, increase a transmission speed and reduce the employed number of memories is shown in FIG. 2. In FIG. 2, one DRAM 42 is connected with first and second processors 10 and 26 through buses B1 and B2, in contrast to the system of FIG. 1. In the structure of multi processor system of FIG. 2, to access to one DRAM 42 by both processors through two respective paths, two ports are required on the DRAM to be connected corresponding to the buses B1 and B2.\nHowever, a DRAM according to a conventional art is a memory 1 having a single port PO as shown in FIG. 3. FIG. 3 illustrating a conventional structure of DRAM, a memory cell array includes first to fourth banks 3, 4, 5 and 6 each having a row decoder and a column decoder. An upper input/output sense amplifier and driver 13 is operatively connected to first bank 3 or third bank 5 through a multiplexer 11, 12. A lower input/output sense amplifier and driver 15 is operatively connected to a second bank 4 or fourth bank 6 through a multiplexer 13 or 28.\nFor example, when a memory cell within the first bank 3 is selected and data stored in the selected memory cell is read, an output operation of the read data is described as follows. First, a selected word line is activated, then data of a memory cell sensed and amplified by a bit line sense amplifier is transferred to a local input/output line pair 9 by an activation of corresponding column selection line CSL. Data transferred to the local input/output line pair 9 is transferred to a global input/output line pair 10 by a switching operation of first multiplexer 21. The second multiplexer 11, connected to a global input/output line pair 10, transfers data of the global input/output line pair 10 to the upper input/output sense amplifier and driver 13. Data sensed and amplified by the input/output sense amplifier and driver 13 is output to a data output line L5 through a path part 16.\nHowever, when data stored in a memory cell of the fourth bank 6 is read, the data is output to an output terminal DQ, sequentially passing through multiplexer 24, multiplexer 28, lower input/output sense amplifier and driver 15, path part 16 and data output line L5. That is, DRAM of FIG. 3 has the structure that two banks share an input/output sense amplifier and driver, and is a single port memory in which the input/output of data is performed through one port PO. The DRAM 1 of FIG. 3 can be used in the system of FIG. 1, but it may be difficult or impossible to use in the multi processor system referred to in FIG. 2.\nReferring to FIG. 4, illustrating a multi processor system 50, a memory array 35 is constructed of first, second and third portions. The first portion 33 of the memory array 35 is accessible only by a first processor 70 through a port 37. The second portion 31 is accessible only by a second processor 80 through a port 38. The third portion 32 is accessible by both of the first and second processors 70 and 80. The size of first and second portions 33 and 31 of the memory array 35 can be variously changed according to an operating load of the first and second processors 70 and 80. Furthermore, the memory array 35 may be a memory type or disk storage type.\nTo realize third portion 32 shared by the first and second processors 70 and 80 within the memory array 35 in the structure of DRAM, several difficulties arise. For example, a layout of memory regions and input/output sense amplifiers within the memory array 35, and a proper read/write path control technique require attention. In the read/write operation through respective ports an operating time interval must be shortened to obtain efficient port use permissions and a high data transmission speed.\nAs a result, there remains a need for an improved method and apparatus for sharing a shared memory region allocated within a DRAM memory cell array in a multi processor system having two or more processors, and to enable transmission of data between ports."}
{"text": "1. Field of the Invention\nThis invention generally relates to set screws and, more particularly, to set screws of the \"cup-point\" type.\n2. Description of the Prior Art\nSet screws are used in various applications to removeably secure fly wheels, pinions, pulleys, cams, propellers and the like to shafts. In such applications, the set screws experience both vibrations and shock loads which might tend to loosen them.\nMany set screw points have been designed in the past to increase the ability of the set screw to resist vibration. Examples of such set screws include: U.S. Pat. No. 1,330,792, Frick; U.S. Pat. No. 2,235,435, Koester; U.S. Pat. No. 2,245,525, Dicely; U.S. Pat. No. 2,462,910, Simmons; U.S. Pat. No. 2,778,265, Brown; U.S. Pat. No. 2,295,314, Whitney; U.S. Pat. No. 2,907,245, Whitson; U.S. Pat. No. 3,286,754, Klooz et al.; U.S. Pat. No. 3,419,058, Walker; and U.S. Pat. No. 3,920,060, Heldmann et al.\nOf particular interest is U.S. Pat. No. 2,314,274, Hallowell, which discloses certain types of set screws having knurls on the work-engaging end wherein the work-engaging extremity is provided with centering means that will be effective to establish the position of the screw with the work prior to engagement of the interlocking means. Also, U.S. Pat. No. 2,992,669, Fesmire, discloses a certain knurled cup-point set screw wherein the cup is counter-bored Additionally, U.S. Pat. No. 4,764,068, Crispell, discloses a set screw having a threaded shank and a cup-point formed at one end of the shank wherein the cup-point has an outer frustoconical surface and an inner cylindrical surface which together define an annular wedge terminating in a circular edge. A portion of the outer frustoconical surface remote from the circular edge is provided with knurling.\nMost of the above-mentioned prior art set screws have cup-point configurations formed by a machining process. However, none of these prior art references suggest the novel work-engaging end portion of the set screw of the present invention which can be formed by a forging process."}
{"text": "It is well known to those skilled in the art that intraocular lenses are predominantly designed to replace a previously or simultaneously removed lens in a cataract patient (see, for example, U.S. Pat. No. 5,628,798). However, although the implantation of intraocular lenses has constituted an appreciable surgical advance, such implantation has been known to cause immediate or latent damage to the corneal endothelium, immediate or latent inflammatory responses in the anterior and/or posterior segments of the eye, immediate or latent secondary fibrosis and/or neovascularization, and other problems.\nIntraocular lenses have been surgically implanted into an aphakic eye in order to take the place of the natural lens. (See, U.S. Pat. No. 2,834,023). Intraocular lenses have been implanted into the posterior chamber of the phakic eye, i.e., an eye having a natural lens in situ, to compensate for refractive errors or to create a specific refraction to assist in visual function. (See, U.S. Pat. No. 4,769,035).\nCertain difficulties are associated with implanting an intraocular lens in the phakic eye that are not encountered when implanting a lens in the aphakic eye. The phakic eye is a substantially more reactive environment than the aphakic eye. Inflammatory reactions tend to be greater in the phakic eye resulting in a concomitant increase in damage to the eye caused by implanting intraocular lenses. One reason is that in the aphakic eye, the natural lens does not pull on the highly reactive ciliary body thus, the ciliary body is in a \"resting state\" and tends to undergo some degree of atrophy. Additionally, the presence of the natural lens in the phakic eye crowds the area in which an intraocular implant can be placed in the eye.\nPlacements of intraocular lenses in the posterior chamber of the phakic eye also have been known to cause cataract formation in the natural lens that remains in situ due to contact between the implant and the natural lens. In contrast, implanting intraocular lenses in patients having cataract removal cannot induce such an effect since the natural lens has been removed.\nThe anterior chamber of an eye is that area in front of the iris and behind the cornea. The iris acts as a divider between the anterior chamber and the posterior chamber. The anterior chamber was originally studied as a preferred location for aphakic intraocular implants particularly when no posterior capsule was present. However, significant drawbacks were discovered.\nThe phakic eye has a shallower anterior chamber (i.e., the average antero-posterior depth is less) than the aphakic eye, and the iris is in contact with the anatomic lens. Therefore, if an inflammatory reaction occurs in the phakic eye, there can be adherence of the iris to the anterior surface of the anatomic lens. Furthermore, if posterior chamber implants were to cause an inflammatory reaction in the phakic eye, cataract formation may occur in the natural lens which remains in situ.\nFor intraocular lenses implantation, it was known in the art that the preferred location of the implant was in the posterior chamber, i.e., that area behind the iris and in front of the natural lens in situ. Intraocular implants for the posterior chamber have been designed to treat myopia (nearsightedness) and hyperopia (farsightedness).\nHowever, a drawback existed for correcting refractive errors in patients with high myopia, extreme nearsightedness, with posterior chamber ocular lenses. Because of the high degree of corrective refraction of the optic, which causes the optical zone to be quite small, a myopic posterior chamber lens may irritate the iris and potentially touch the natural lens. Anterior chamber intraocular lenses for high myopia were explored for implantation in the phakic eye which would not be in contact with uveal tissue or the human lens. Hyperopic posterior lenses would not cause similar damage because the edge of a hyperopic lens is thinner and does not project the peripheral iris forward. In addition, the vault of a hyperopic posterior lens prevents human lens touch. Therefore, there is no incentive to place a hyperopic lens into the anterior chamber.\nU.S. Pat. No. 4,676,792 discloses an uncoated minus power anterior chamber implant for the treatment of myopia. Uncoated minus power implants have been used in patients having high myopia that is not satisfactorily corrected with spectacles or contact lenses (see, for example, Colin et al., 1990, Refractive and Corneal Surgery 6:245-251 and Baikoffet al., 1990, Refractive and Corneal Surgery 6:252-260). The prior art does not disclose implanting a positive power implant into the anterior chamber of a phakic eye for the treatment of hyperopia. Moreover, due to the shallowness of the anterior chamber in hyperopic eyes, surgical location of a positive power implant in the anterior chamber of the phakic eye for the treatment of hyperopia would be in contravention of, and counterintuitive to, the state of the art.\nIt has been generally acknowledged by those skilled in the art that there are significant risks involved with the use of anterior chamber angle supported implants in the aphakic eye (reviewed in Apple et al., in Intraocular Lenses: Evolution, Designs, Complications and Pathology (William and Wilkens, Baltimore) 1989, Chapter 4, pp. 59-105) and more so in the phakic eye (Id. at p. 65, col. 1). For example, when such an implant is inserted into the eye, temporary or permanent adhesions of the implant to iris tissue may result, causing damage to these structures to ensue either immediately or over the long term affecting pupillary mobility and contour. In addition, once the implant is in position, it may cause similar angle adhesions due to mechanical and/or chemical inflammation which may lead to fibrosis of a progressive nature. This would make subsequent removal of the implant a complex, dangerous surgical procedure. Other problems associated with such implants are cataract formation, secondary glaucoma, corneal edema, hyphema, and progressive endothelial cell loss, in addition to other complications.\nAs observed in Ophthalmology Alert, Vol. 1, No. 11 (Nov. 1990), pp. 41-42, Comment on page 42, several American manufacturing companies that were preparing to begin clinical trials of phakic anterior chamber ocular implants in the United States are now likely to abandon these studies, due to the attendant risks associated with anterior chamber implants in the phakic eye and the difficulty of obtaining approval of the U.S. Food and Drug Administration (FDA) for the use of the implants. A significant risk involved in the use of such anterior chamber implants in the phakic eye is the potential for the implanted lens to contact the corneal endothelium, the anatomic lens or the iris with resultant complications, even with the enhanced vault design theorized in the studies discussed in Ophthalmology Alert because such design would bring the lens optic edge quite close to the midperipheral corneal endothelium. One alternative to avoid contact with the corneal endothelium is to reduce the diameter of the optic of the minus power lens, however, such a modification creates significant drawbacks, including glare and haloing under low light conditions, ovalization of the pupil, and, as reported in Perez-Santonja et al., J. Cataract Refract Surg 22:183-187 (1996) (discussing Baikoff ZB5M), induces an inflammatory response.\nIt would clearly be advantageous to employ an uncoated positive power anterior chamber ocular implant that would avoid the occurrence of the above-described problems associated with anterior chamber implants for the treatment of hyperopia.\nHyperopia, a condition commonly known as farsightedness, has been known to be treated by wearing positive power eyeglasses, contact lenses, and/or surgically implanted iris supported or posterior chamber intraocular lenses. Plastic rings have also been surgically sutured to the sclera for the treatment of hyperopia and presbyopia. It was generally believed by persons of skill in the art that the anterior chamber of the phakic eye of those individuals suffering from hyperopia was more shallow than in the phakic eye of those individuals suffering from myopia. Therefore, attempts to correct hyperopia did not involve implantation of a positive power implant into the anterior chamber of the phakic eye.\nU.S. Pat. No. 4,769,035 discloses a posterior chamber implant for implantation in the phakic eye for correction of hyperopia. This posterior chamber implant is designed specifically to come into immediate and direct contact with the anterior surface of the natural lens. The surgical procedure for implanting this posterior chamber implant requires a pre-implantation determination of the precise shape and contouring of the surface of the natural lens and the specific amount of optical correction required for placement of such a lens in direct contact with the natural lens.\nU.S. Pat. No. 5,489,299 discloses a method of treating hyperopia and presbyopia wherein the effective working distance of the ciliary muscle can be increased by expanding the sclera in the region of the ciliary body by suturing to the sclera a relatively rigid band having a diameter slightly greater than that of the sclera in the region of the ciliary body. Rings are surgically implanted, and intraocular lenses are not used at all in this procedure.\nU.S. Pat. No. 5,192,319 discloses an intraocular lens surgically positioned in the anterior chamber to be used in addition to the natural lens to correct refraction by fixation of such lens on the iris using one or more pairs of pincer-like extensions for holding the position of the lens to the iris tissue. This patent discloses a positive power implant for placement in the anterior chamber of the phakic eye for correction of hyperopia, specifically designed to attach to the iris. Although such positioning on the iris may provide for safe clearance away from vital structures such as the natural lens and the corneal endothelium, the implant has a number of significant drawbacks. Damage may occur to the sensitive iris tissue and the damage is increased when adjusting the position of the implant in cooperation with the existing lens. Also, deterioration of the iris tissue occurs, thereby decreasing the ability to place future replacement implants on the iris. Contact with vital intraocular structures may occur. The patent states at column 4, lines 25-27, that \"[a]ll designs are a compromise between the anatomical limitations caused by the position of the intraocular lens in the anterior chamber and the optical requirements.\" In addition the disadvantages of iris mounted anterior chamber implants are discussed by Apple et al., in Intraocular Lenses: Evolution, Designs, Complications and Pathology (William and Wilkens, Baltimore) 1989 Chapter 3, pp. 45-58, which states on page 58 that such lenses have become obsolete due to numerous complications.\nU.S. Pat. No. 4,871,363 also discloses a corrective intraocular lens positioned in the anterior chamber in a phakic eye for correcting high myopia wherein a leading haptic is shorter than the trailing haptic. Such a lens would necessarily be a minus power lens, even though the illustration of FIG. 3 is drawn showing what appears to be a positive power lens, because it is designed specifically for treatment of myopia. The specification of '363 does not provide any disclosure of a positive power ocular lens in a phakic eye for treatment of hyperopia since the invention is limited to the treatment of myopia. Moreover, the intraocular lens that is disclosed in the '363 patent is admitted to be designed primarily for the replacement of a natural lens in an aphakic eye at col. 3, lines 36-44. The '363 patent does not appreciate the problems associated with positive power ocular implants actually implanted into the phakic eye.\nU.S. Pat. No. 4,676,792 discloses a minus power intraocular implant for placement in the anterior chamber of the phakic eye for the treatment of myopia. Although FIG. 2 of the '792 patent shows a convex anterior surface of the optic, the optic is minus power due to the concave curvature of the posterior surface of the optic. The '792 patent does not appreciate the complications associated with a planar or convex posterior surface of the optic in combination with a convex or planar anterior surface.\nTo Applicant's knowledge, there has not been a successful implantation of a positive power anterior chamber ocular implant angle supported in the phakic eye for the treatment of hyperopia. In view of the problems and risks encountered with the development of minus power anterior chamber intraocular lenses, it would have been unexpected for positive power anterior chamber intraocular lenses to be successfully modified for use in the phakic eye for the treatment of hyperopia.\nIt is an object of the present invention to provide a positive power ocular implant for placement in the anterior chamber of a phakic eye to correct refractive errors caused by hyperopia."}
{"text": "Field of the Invention\nThe present invention relates to a position detection apparatus that detects a position of a target object.\nDescription of the Related Art\nThere is a known image forming apparatus that primarily transfers toner images respectively formed on a plurality of photosensitive members to an intermediate transfer belt and secondarily transfers a color image composited on the intermediate transfer belt to a recording material.\nIncidentally, when the intermediate transfer belt in the image forming apparatus is deviated in a width direction that intersects perpendicularly with a belt conveying direction, color misregistration in which the toner images of the plurality of colors on the intermediate transfer belt are deviated may occur. In order to prevent generating such color misregistration, there is a known correction control that detects a deviation amount in the width direction that intersects perpendicularly with the belt conveying direction of the intermediate transfer belt and that corrects belt driving corresponding to the deviation amount.\nAs an apparatus that detects a deviation amount of an intermediate transfer belt, there is a proposed apparatus that is provided with a swinging arm that is in contact with an edge of the intermediate transfer belt and swings and two optical sensors that are arranged on the swinging arm so as to shift in a longitudinal direction thereof that is the conveyance direction of the intermediate transfer belt. In this belt-deviation-amount detection apparatus, the swinging arm swings corresponding to the deviation amount of the intermediate transfer belt. And a result of whether the deviation amount of the belt falls in tolerance is detected using changes of the signals from the two optical sensors on the swinging arm due to swinging (for example, see Japanese Laid-Open Patent Publication (Kokai) No. 2010-243791 (JP 2010-243791A)). Moreover, this belt-deviation-amount detection apparatus detects the deviation amount of the belt in five levels.\nHowever, the above-mentioned conventional apparatus is easily affected by variations of the sensor positions and skew of a target object (belt), and particularly, there is a problem of erroneous detection near a boundary between adjacent deviation levels."}
{"text": "1. The Field of the Invention\nThis invention relates generally to the field of wireless communication software. In particular, embodiments of the present invention relate to a method for locating and connecting to a wireless device within a wireless network.\n2. The Relevant Technology\nShort range wireless communication schemes are increasingly used as an alternative to wire-based networks to interconnect computers and computer devices. The use of wireless communications provides a number of advantages, including the elimination of complex, expensive, and inconvenient wire-based connection schemes. While a number of wireless communications standards have been developed for such applications, one popular example is known as “Bluetooth.” The Bluetooth standard uses low power radio frequencies to allow communication between various devices such as mobile phones, laptop and desktop computers, printers, modems, PDAs, etc. This technology was originally envisioned for the purpose of replacing cabling and other hard-wired connection schemes used to connect auxiliary devices to a desktop or laptop computer. However, Bluetooth has further evolved into a method of sending both data and voice signals between a wide range of devices. For example, a Bluetooth-enabled PDA can be configured to automatically connect to a Bluetooth-enabled communications line within a building, therefore gaining wireless access to computing resources, such as printers, an Internet portal, etc.\nIn an effort to ensure compatibility among Bluetooth implementations, a Bluetooth special interest group (SIG) was formed. One of the main objectives of the Bluetooth SIG is to formulate a standardized specification for every element of Bluetooth. The specification dictates numerous software and hardware characteristics so that devices made by different manufacturers will be compatible with one another. The specification dictates all aspects of Bluetooth technology from the software stack necessary to run Bluetooth devices to the actual hardware requirements.\nAs per this standard, Bluetooth compatible technology is composed of multiple components that together provide for wireless communication capabilities. For example, each Bluetooth-compatible device is configured with a radio that can transmit and receive wireless signals on a particular frequency band. In addition, there are certain structural elements required for connecting the radio to the specific host device. Further, each Bluetooth-compatible device is also configured with particular software to allow the device to properly process and route data or voice signals. This software is typically referred to as the “protocol stack” because it is composed of numerous components that are each configured to perform a particular function. Information is generally passed between the respective layers of the stack either from the bottom of the stack to the top or vice versa, depending on the direction of the communication involved. The protocol stack is further broken up into a “transport protocol group,” a “middleware protocol group” and an “application protocol group”.\nIn general, the transport protocol group includes functional components that are designed to allow Bluetooth devices to locate one another within a Bluetooth network, and to manage the links that allow higher layer protocols and applications to pass data through the transport protocols. The transport protocol group is comprised of a radio, a baseband controller, a link manager, a logical link control and adaptation protocol (L2CAP), and a host controller interface (HCI). The baseband controller defines how devices search for one another and how they connect to one another via a Bluetooth Device Discovery and Connection process. The baseband controller also defines the master and slave roles between devices. The link manager supervises the creation of pairing between devices and the encryption of data. The link manager also supervises power control among devices. For example, the link manager can initiate a low power adjustment between devices to conserve energy. The L2CAP supports the protocol multiplexing that allows multiple protocols to utilize the same air-interface. It also maintains the desired level of service among devices. The HCI allows higher layers of the stack to access the baseband controller and link manager. The HCI is not a required part of the Bluetooth specification. The radio, link manager and baseband controller are sometimes grouped together to form what is referred to as a Bluetooth module. The remaining components of the transport layer are typically located on the actual Bluetooth host device.\nThe middleware protocol group is directed to functionality that facilitates communication between host applications and the Bluetooth module. The protocols can include both third-party protocols and industry standard protocols specified by the Bluetooth SIG. The middleware protocol group specifically includes a RFCOMM, a service discovery protocol (SDP), a set of IrDa interoperability protocols, and a telephony control protocol (TCS). The RFCOMM protocol creates virtual serial ports to facilitate the replacement of cables with a wireless Bluetooth interface. The SDP protocol defines a standard method for Bluetooth devices to discover and learn about the services offered by other Bluetooth devices within the network. The SDP protocol is often utilized by other protocols for the purpose of identifying what services are available over a wireless air interface. The set of IrDa interoperability protocols helps to define the syntax of the data to be exchanged between devices. The TCS protocol is designed to control the exchange of voice communications between devices.\nThe application protocol group consists of the various applications that utilize Bluetooth communications. The applications may be unaware that data is being transmitted over a Bluetooth air-interface rather than a standard cable, or the application may be specifically designed to interact with Bluetooth devices.\nBluetooth devices from what are known as “piconets” with one another. A piconet is comprised of one Bluetooth device acting as a master and numerous other devices acting as slaves. Both the master and the slave devices can transmit and receive data from one another. However, the master device defines the data hopping sequence and is usually the device that initiates the communication. Initiation of a communication depends, in part, on the configuration state of a device. For example, an “active state” is a system configuration in which the slave device is always listening for transmissions from the master. In a “parked state” the slave device is only periodically listening to the master. Thus, a parked slave must become active before it can communicate normally with the master. In addition, there is a sniff state and a hold state, both of which are considered active states. The sniff and hold states are used as alternatives to the full active state for the purpose of conserving power. A slave device in a sniff mode essentially communicates with the master at regular intervals rather than always listening. A slave device in a hold mode essentially stops listening to the master for a specified period of time.\nThere are several limitations present in the Bluetooth standard that can degrade the overall operating efficiency of a Bluetooth network. A primary limitation relates to the number of active slave devices that can be connected within a single piconet. Up to seven active slave devices and up to 255 parked slave devices can be linked to a master device in a single piconet at a given time. The seven device limitation prevents the master from actively communicating with more than seven slave devices at any given time. This device limitations is mandated by the physical bandwith characteristics of Bluetooth devices. Obviously, the seven device limitations can limit the functionality that is available to a master device at any given time, which may in turn severely limit the operating efficiency of the Bluetooth network. For example, if an application needs to establish a communication link with the Internet, it may request a Bluetooth-connection with a Bluetooth enabled modem. However, if the modem is attached to one of the parked Bluetooth slave devices (i.e., not available as one of the seven active devices), it may take a relatively long amount of time before the parked device becomes active and actually processes the request from the application. This time delay is obviously undesirable.\nIn addition, one Bluetooth device may be part of multiple piconets. For example, the master of one piconet may also be a slave in another piconet. In this scenario, a slave device from one piconet could order a “print” command. That print command could be routed through the master to a slave which happens to be a master of another piconet that contains a printer. That second master will then send the signal to the printer. However, this type of situation could severely limit the bandwidth of the second master as it is routing the print command (and associated print data) to the printer, thereby limiting its ability to respond to other requests within its piconet. Again, this problem is in large part a result of the seven device limitation.\nIn addition to communicating with other devices, each individual Bluetooth device periodically executes a search routine (SDP protocol) to locate new devices within the area. This is necessary for automatically connecting to new devices as they become available. In particular, the master periodically searches for new devices in addition to communicating with all of the slaves in the particular piconet. Although it is very useful for the master of each piconet to periodically find new stationary devices as they become available, it is not always possible to connect with these devices in a timely manner. For example, a requesting device connected to the computer network may desire a connection with another device that is known but not connected to the network. An automatic SDP search routine that is executed every ten minutes may find the device but it would most likely introduce a long delay. Alternatively, it would be possible to increase the automatic SDP search routine but this would introduce delay in all of the other transmissions to and from the master.\nTherefore, a method for efficiently connecting with a specific wireless device at an unknown location within a wireless computer network is needed. In addition, the method should be cost effective to implement in a wireless networking environment. The method should require less processing power and bandwidth than the standard SDP search routine."}
{"text": "Caps and hats vary greatly in terms of design and functionality. While baseball hats are very trendy in contemporary times, they fail to provide the warmth necessary for extreme drops in temperature. A typical baseball cap is constructed of a webbed plastic mesh. Such a hat prevents sun from obstructing the wearer's view, but the hat does little in the way of maintaining body heat. In fact, baseball hats are generally worn as ornamental pieces outside the realm of sporting events. Although some baseball hats are constructed of heavy cotton and other materials, the basic skeleton of a baseball hat forces the hat to sit atop the wearer's head, not on and around the wearer's head.\nSnow hats, on the other hand, are traditionally constructed primarily of wool and acrylic. They are designed to keep the wearer's head warm. In this regard, a typical snow hat exhibits uniform construction and has little structural integrity when removed from the wearer's head. Snow hats are not intended to be rigged and present obstructions to the sun, wind, and rain; but rather, snow hats are designed to be close fitting and insulating with respect to the wearer's head, maintaining anterior temperature as moisture and heat are conserved.\nThe prior art is replete with various types of hats, none of which approach the design and functionality of the present invention.\nU.S. Pat. No. 80,352, issued to Ibach and Weidenman on Jul. 28, 1868, illustrates a method of attaching a paper visor to a paper cap wherein a visor and a loop of material is fitted to the bottom rim of a skullcap. When the cap becomes dirty, the visor can be attached to another cap, whereas in the present invention, the visor must firmly remain in combination with a head covering to remain fully effective. Moreover, unlike the present invention, Ibach and Weidenman's device does not flexibly surround the user's head to promote heat retention. U.S. Pat. No. 442,921, issued to Stohr on Dec. 16, 1890, shows a knitted cap with a visor in which a hat with depending ear flap and ties is combined with a visor. Stohr, unlike the present invention, does not provide any means of capturing moisture from the brow or forehead of the user during athletic activity. Furthermore, Stohr's invention, unlike the present invention, employs incisions in the cap to attach the visor.\nU.S. Pat. No. 2,149,655, issued to Yamaguchi on Aug. 3, 1938, illustrates a head covering made entirely of crocheting yarn. The device discloses a means for reinforcing a crocheted visor to a crocheted hat. Unlike the present invention, Yamaguchi's device is concerned with forming a bead at the junction of the crown and the visor to add stability to stiching.\nU.S. Pat. No. 2,158,861, issued to Meyer on Sep. 2, 1937, shows a visor for collapsible caps which is capable of assuming a flat position when a collapsible cap is in a folded, flat position. Unlike the present invention, Meyer's device is not capable of providing a visor in combination with a knit type hat.\nU.S. Pat. No. 2,417,986, issued to Marder et al. on Mar. 25, 1947, is directed to a cap and visor with draw strings extending to the nape of the user's neck. Unlike the present invention, Marder's invention has multiple lines of stitching and is secured behind the user's head.\nU.S. Pat. No. 2,420,569, issued to Sewell on May 13, 1947, shows a baseball cap. Unlike the present invention, Sewell's device has a visor which is attached to the crown of the cap.\nU.S. Pat. No. 2,651,044, issued to Stankiewicz et al. on Sep. 8, 1953, depicts a scarf hat having interfitting concentric outer and inner head band members for receiving between them an edge of a scarf. Unlike the present invention, the visor is attached by fitting between the outer and inner head band members.\nU.S. Pat. No. 4,601,070, issued to Sargentini on Jul. 22, 1986, shows a novelty ski hat. Unlike the present invention, there is no visor attached to the hat.\nU.S. Pat. No. 5,471,684, issued to Casale on Dec. 5, 1995, shows a sports cap with a sliding brim. Unlike the present invention, the brim is detachably secured to the bottom edge of the cap. Furthermore, unlike the present invention, the majority of the cap can be detached from the bottom edge of the cap.\nU.S. Pat. No. 5,481,759, issued to Rinaldi on Jan. 9, 1996, illustrates an expandable baseball hat and cover. Unlike the present invention, Rinaldi's device employs tabs engaged to gather flexible material for a snug fit. The present invention, by nature of its design, fits snugly to the user's head without employing a complicated tab system.\nFrench Publication No. 621,305, invented by Menant and published on May 9, 1927, shows a woven cap. Unlike the present invention, Menant's device employs multiple stiches about the middle periphery of the cap. Further, Menant's device is not designed, by virtue of its shape and inability to snugly fit the user's head, for athletic use.\nAccordingly, the need arises for a multi functional hat with a brim or visor, which is a blends the functionality of both the typical snow hat and the typical baseball cap. The multi functional hat must be rigid so as to protect the wearer's face from the outdoor elements. Moreover, the multi functional hat should be capable of maintaining the warmth of the wearer's head despite baseball cap type adaptations. The multi functional hat should be of a design that does not compromise the advantages of existing snow hats, while at the same time, exploiting the current weaknesses of headgear design."}
{"text": "1. Field of the Invention\nThe present invention relates to a thermal management system, and a method for managing thermal characteristics, for a heat producing system.\n2. Background Art\nIn response to demands for improved fuel economy and reduced emissions, vehicles today are being manufactured with systems designed to increase combustion efficiency and reduce parasitic losses of various vehicle components. One way to increase combustion efficiency in an internal combustion engine is to maintain a high degree of control over the temperature of the combustion in the engine cylinders. The use of an effective vehicle thermal management system can help to achieve this goal. For example, controlling one or more of the engine oil temperature, the engine coolant temperature, and the intake air temperature, can provide an effective means for ensuring that combustion within the engine cylinders takes place within a desired temperature range. Controlling the temperature of the combustion within the engine can help to increase combustion efficiency, and reduce exhaust emissions.\nA number of thermal management systems are described in a Society of Automotive Engineers (SAE) Technical Paper, Document Number 2001-01-1732, entitled “Thermal Management Evolution and Controlled Coolant Flow,” copyright 2001. One such system includes a controllable electric pump for circulating engine coolant through an EGR cooler. The electric pump can replace a larger, mechanical pump, thereby providing an overall space savings. Another system described in the SAE paper includes a separate EGR cooling loop having its own coolant loop separate from the engine coolant loop. The EGR cooling loop includes a controllable electric pump, and its own liquid-to-air heat exchanger for dissipating heat from the EGR coolant.\nWhile a vehicle thermal management system can be used to control the temperatures of various vehicle systems, including the temperature of combustion, it would be desirable if the same thermal management system could be used to decrease parasitic losses of various components within the vehicle. For example, a thermal management system may employ the use of one or more electric fluid pumps, electric valves and electric fans. These electric components may replace one or more mechanical components which typically operate in accordance with the speed of the engine. Through the use of electric components, controlled by an electronic controller, it would be desirable if such a thermal management system could optimize the operation of the components to reduce overall power consumption while still providing the functionality necessary for an efficient thermal management system."}
{"text": "Conventional database servers receive transactions from user terminals. Upon receipt of a transaction, the database server carries out various processes; for example, it adds, reads out, updates, and deletes data. If this type of database server concurrently receives transactions that attempt to update the same data from a plurality of users, the database server performs exclusive control, such as a lock method etc., on various processes to assure the integrity of the data.\nSpecifically, to assure data integrity, the database server exclusively controls transactions that attempt to update the same data with conflicting transactions. For example, the database server uses pessimistic lock, optimistic lock, or another type of exclusive control etc.\nPessimistic lock is a method that assumes that transactions attempting to update data in a database are highly likely to be received at the same time from a plurality of users. Once a transaction reads out data to be updated, the other transactions are excluded.\nOptimistic lock is a method that assumes that transactions attempting to update data in a database are less likely to be received at the same time from a plurality of users. Once a transaction reads out data to be updated, if another transaction attempts to update and commit the data before the transaction completes the update of the data, the other transaction is treated to be invalid.\nAs described above, to assure data integrity, a database server uses pessimistic lock, optimistic lock, or another type of exclusive control to process conflicting transactions in succession while excluding the other transactions.\nTechnology that enables a database server to process a plurality of transactions together is also disclosed. For example, the database server combines transactions entered in a certain period of time and recombines the transactions for each data to be read out and updated by transactions to determine their execution order, before the database server executes the transactions.\nHowever, the above conventional technology is problematic in that conflicting transactions cannot be executed at the same time. For example, a database server that uses exclusive control executes conflicting transactions in succession to assure data integrity, so the database server cannot execute the conflicting transactions at the same time.\nIn the method in which a plurality of transactions are processed together, if the database server executes a transaction that determines a branch instruction, for example, the database server determines a condition from a result obtained by searching the transaction for data before executing a next process, so the database server cannot send all processes executed in the transaction to a database in advance. It is difficult and is not practical to consider, during the design of an application, updates carried out by other transaction to be executed at the same time and design execution of the transactions.\nJapanese Laid-open Patent Publication No. 2009-271665 is an example of related art."}
{"text": "Oriented polyolefin tape yarns are commonly used in woven fabrics for applications such as carpet backing, geotextiles for soil stabilization and separation, and in coated fabrics such as housewrap, lumber wrap, and laminates for flexible packaging applications. In all these applications, the foamed, oriented tapes impart high strength and stiffness to the fabrics, and depending on the application, other essential properties. Foamed oriented tape yarns are also used as wire insulation and as crack arrestors in concrete.\nCarpets\nIn the past, jute cloth was used as carpet backing for tufted carpets. More recently, woven fabrics of polypropylene tape yarns have been used for carpet backings.\nTufted carpets can be produced by a two-step process involving 1) stitching pile yarns made from synthetic or natural fibers through the primary backing then 2) fixing the stitches on the underside of the primary backing with a latex adhesive.\nTufted primary backings may have pile yarns that are loops, cut pile, or combinations of cut and loop pile. During the carpet manufacturing process, a second fabric, called the secondary backing may also be used. The latter fabric is laminated to the opposite side of the adhesive layer and contacts the floor. The secondary backing adds dimensional stability to the finished carpet. Typically the secondary backing is a woven fabric with tape yarns in the warp direction and spun yarns in the cross machine (or filling) direction.\nFoamed Tapes for Carpet Backings\nIt is well known in the conventional art that polyolefin resins can be used as the stock materials for forming tape yarns for woven carpet backings. Specifically, it has been found that foamed, oriented polypropylene tape yarns can be used as the stock materials for woven backings. Such foamed tape yarns have ridged surfaces and internal voids. They typically contain less material than ordinary polypropylene tape yarns which are not formed by foaming methods. In addition to weighing less than conventional (unfoamed), oriented tape yarns, foamed, oriented tape yarns for primary carpet backing impart improved tuftability and surface appearance to tufted carpets. In secondary backings, foamed oriented tape yarns impart similar levels of dimensional stability to the finished carpet with less material than conventional (unfoamed) tape yarns.\nGeotextiles\nMany geotextiles are manufactured from oriented polyolefin tapes, including woven fabrics for separation and stabilization of soil layers in roadbeds and for containing soil particles (e.g., silt) at construction sites. Polyolefin tapes are also used in retaining wall fabrics, turf reinforcement mats, ground cover fabrics, geotubes, and erosion control blankets. Polypropylene tapes are used in many of these applications.\nFoamed Tapes for Geotextiles\nWoven geotextiles require high strength and stiffness, and specific levels of permeability to water. Often the strength of the starting tapes is diminished when they are woven into fabrics. In general, foamed, oriented tapes lose less strength in the weaving process than conventional (unfoamed) tapes, resulting in increased fabric strengths compared to fabrics woven from unfoamed tapes with the same starting tensile properties. For applications that require woven geotextiles to act as a filter, such as silt fence, foamed tapes generally have greater dimensions than conventional (unfoamed) tapes of similar weight, allowing less fabric to create the same opening size in the filter membrane.\nPackaging\nOriented polyolefin tapes are commonly used to make reinforcing scrims in flexible packaging applications, such as lumber wrap, bale packaging for cotton and synthetic staple fiber, and retail packaging of food, seeds, and other granular material. The woven scrims are often extrusion coated with a layer of polyolefin resin to create a moisture-proof barrier and increase puncture resistance, strength and stiffness.\nFoamed Oriented Polyolefin Tapes in Packaging\nWoven fabrics for packaging often have critical requirements such as high coverage and strong adhesion of the extrusion coating to the woven scrim. Coverage, which refers to the extent of open area in the woven fabric, depends on the widths of the tapes and count of the tapes in the warp and fill direction. Compared to woven scrims made with conventional (unfoamed) oriented tapes, woven scrims with foamed, oriented tapes can have greater coverage for equal weight (or alternatively, equal coverage with less weight) and improved adhesion. As described in more detail below, improved adhesion is a consequence of the special ridged surface structure of the foamed, oriented tapes made by the method and system of this invention.\nHousewrap\nHousewrap is a moisture resistant, breathable fabric applied to residential and light commercial buildings prior to the installation of the exterior finish, such as shingles, brick, or siding. One purpose of housewrap is to conserve energy by reducing air infiltration. Many housewraps consist of extrusion coated woven scrims made from oriented, polyolefin tapes. They are similar to extrusion coated scrims for packaging, but often have less coverage (more open area in the scrim). These coated fabrics are often mechanically perforated to allow water vapor to pass through. Most housewraps made with woven polyolefin reinforcement easily pass the minimum requirements for tensile strength required by industry regulations.\nFoamed Oriented Polyolefin Tapes in Housewrap\nFoamed, oriented tapes have sufficient strength to produce housewraps that meet industry standards. They also have a ridged surface that imparts improved adhesion of the polyolefin coating to the tapes, another important property for these products. Scrims made from foamed, oriented tapes require less resin than scrims made from similar conventional (unfoamed) tapes while still delivering fully satisfactory tensile properties to the final product.\nOther Applications\nOriented polyolefin tapes are used for other applications such as netting for produce bags, open knitted fabrics for curing meat, and for wire insulation. Polyolefin tapes that are cut into 1 to 2-inch lengths are also used as crack arrestors in concrete mixtures.\nFoamed Oriented Polyolefin Tapes in Other Applications\nFoamed, oriented tapes for other applications, can include, but are not limited to netting applications and open knitted fabrics because foamed tapes impart sufficient strength and stiffness at lower weight than conventional, unfoamed tapes. When used as wire wrap, foamed, oriented tapes impart improved insulation properties because of their increased thickness compared to conventional (unfoamed), oriented tapes. As crack arrestors in concrete, chopped foamed tapes made by the system and process of this invention have surface characteristics that lead to improved mechanical bonding between the foamed, chopped tape and cement mixture compared to an unfoamed, chopped tape of similar length.\nLimitations of Conventional Thermoplastic Foaming Processes\nWhile the conventional art describes the advantages of foamed, oriented polypropylene tapes in the above applications, much of the conventional art does not articulate how these foamed tapes can be manufactured. Most of the conventional art only describes the mechanical properties of the end products (i.e., the tapes and woven fabrics) and fails to provide any information on a sequence of specific steps and/or processes for producing such foamed, oriented tapes. Furthermore, the surface characteristics of the foamed oriented tapes are not described.\nSome of the conventional art mentions specific information on a few, isolated steps that can be used in processes for making foamed, oriented polyolefin tapes for use in primary carpet backings. However, the conventional art does not provide a complete set of steps which would allow one of ordinary skill in the art to produce repeatable and accurate results for manufacturing foamed, oriented polyolefin tapes for carpet backings. Specifically, the conventional art does not provide discrete steps that describe how to manufacture foamed, oriented polyolefin tapes economically and efficiently at high line speeds, such as speeds at or above 850 feet per minute. For efficient production, the number of end breaks in the process must be low—preferably less than about 5% per hour and more preferably less than about 2% per hour.\nThe conventional art also mentions the use of certain equipment for making foamed polyolefin profiles, such as fence posts, moldings, and sheets for thermoforming. FIG. 1A shows a high density foam exiting a die 106 and expanding into a shaping device or calibrator 106B. The purpose of the calibrator 106B is to fix the dimensions of the expanding foam profile at it cools. While this is needed for thick section parts, it is very difficult to implement with thin sheets for tapes that are produced at faster extrusion speeds, since the opening (interior height dimension) within the calibrator for thin sheets would be very small, leading to frequent clogging and tearing of the foam sheet.\nIn conventional art, a set of stack smooth rolls (not shown) is sometimes used to control the thickness and smoothness of an extruded foam sheet. Since foaming processes with stacked rolls are typically run at lower extrusion speeds than the tapes of this invention and are designed to create a smooth surface on the cooled foamed sheet, they are also not well suited for making the tapes of this invention, which have a unique surface structure.\nFIG. 1B shows another method for making foamed, thick section profiles. In FIG. 1B, the foamable resin 107 is spun into filaments through a spinneret die 220. The extruded filaments coalesce into a thick section that is further shaped by a calibrator 106C. The calibrator creates a smooth surface on the coalesced filaments. In the method of this invention, a smooth surface is not required or desired. Consequently, the conventional processes of shaping a foamable extrudate, either by a calibrator or stack of smooth calendar rolls do not provide much useful information about how to make the thin, foamed tapes with ridged surfaces that are the end products of the inventive method and system.\nAccordingly, there is a need in the art for a method and system that describes the discrete steps and equipment that can be used to produce foamed, oriented polyolefin tape yarns for the above applications. Specifically, there is a further need in the art for a method and system that describes how such foamed, oriented polyolefin tapes can be manufactured efficiently at line speeds that allow them to be a lower cost alternative to conventional (unfoamed) tapes of similar dimensions."}
{"text": "In accordance with requirement on so-called down-sizing to reduce in size and weight of electronic and electric apparatuses, same kind of performance is required on the printed wiring boards also. For the printed wiring boards, formation of a miniaturized and high functional wiring is required in a limited mounting area. Thus, a fine-pitch wiring must be formed to prepare a high density wiring. To obtain such a fine-pitch wiring, a thinner copper foil having the lower surface roughness on the bonding side to a base material has been employed to reduce over-etching time.\nFor this purpose, a low-profile electro-deposited copper foil is generally used. Also, to ease handling of the copper foil and a copper clad laminate when the copper foil thickness is reduced, efforts have been made to increase mechanical strength of the copper foil. Such a low-profile electro-deposited copper foil having excellent mechanical strength is disclosed in Patent Documents 1 and 2, which will be more specifically described below.\nObject of Patent Document 1 is to provide a low-profile electro-deposited copper foil, which has a low-profile surface suitable for practical use in a printed wiring board and a negative-electrode collector of a lithium ion secondary battery as well as excellent fatigue ductility. More specifically, object of the document is to provide a low-profile electro-deposited copper foil having a surface roughness (Rz) of 2.0 μm or less, which has a low-profile surface without wavy undulation and having an elongation of 10.0% or more at 180° C. For this purpose, Patent Document 1 discloses a method of manufacturing an electro-deposited copper foil by using an aqueous sulfuric acid/copper sulfate solution as an electrolytic solution, an insoluble anode made of a titanium plate coated with a platinum metal or an oxide thereof and a cathode drum (which faces to the anode) made of a titanium and supplying direct current between the electrodes. In this manufacturing method, the electrolytic solution contains an oxyethylene surfactant, polyethylene imine or a derivative thereof, a sulfonic acid salt of an active organic sulfur compound and chloride ions. By the presence of these compounds, a low-profile electro-deposited copper foil (which has a low-profile surface with uniform roughness and having a surface roughness (Rz) of 2.0 μm or less without wavy undulation and having an elongation of 10.0% or more at 180° C.) can be obtained. Examples of Patent Document 1 discloses that the electro-deposited copper foil obtained had a deposit surface having a surface roughness (Rz) of 0.9 μm to 2.0 μm, an elongation as received of 10% to 18%, an elongation at 180° C. of 10% to 20%, a tensile strength as received of 340 MPa to 500 MPa, and a tensile strength at 180° C. of 180 MPa to 280 MPa. The document also discloses that the deposit surface of the electro-deposited copper foil had a gloss [Gs(85°)] of 120 to 132, which was measured at an angle of 85° for the transverse direction.\nPatent Document 2 discloses the technologies which object to provide a low-profile electro-deposited copper foil having a low-profile surface and less tensile strength loss after heating, and having an excellent elongation at high temperature, and object to provide a method for manufacturing the same also. For this purpose, five additives, namely, hydroxyethylcellulose, polyethylene imine, acetylenic glycol, a sulfonic acid salt of an active organic sulfur compound and chloride ions are added in an electrolytic solution, which is an aqueous sulfuric acid/copper sulfate solution. By the presence of these five additives, it enables to obtain a low-profile electro-deposited copper foil having a surface roughness (Rz) of 2.5 μm or less, a tensile strength (at 25° C. measured within 20 minutes after finishing electro-deposition) of 500 MPa or more, a reduction rate of the tensile strength (at 25° C. measured at 300 minutes after finishing electro-deposition) of 10% or less or a tensile strength loss (at 25° C. measured after a foil is heated at 100° C. for 10 minutes after finishing electro-deposition) of 10% or less, and an elongation at 180° C. of 6% or more.\nExamples of Patent Document 2 disclose that\nan electrolytic solution, which consists of an aqueous sulfuric acid/copper sulfate solution containing 100 g/L sulfuric acid (H2SO4) and 280 g/L copper sulfate 5 hydrates (CuSO45H2O), was used as a basic solution;\nadditives, hydroxyethylcellulose, polyethylene imine, sodium 3-mercapto-1-propanesulfonic acid, acetylenic glycol and hydrochloric acid were added to the basic solution;\nthe electrolytic solution was filled in the gap between an insoluble anode made of platinum oxide coated titanium and a cathode drum made of titanium;\nelectrolysis was performed at a current density of 40 A/dm2 and a temperature of the electrolytic solution of 40° C. to obtain an electro-deposited copper foil of 18 μm in thickness; and\nthe electro-deposited copper foil thus obtained had a deposit surface having a surface roughness (Rz) of 1.5 μm to 2.3 μm, a tensile strength as received of 650 MPa to 900 MPa and a tensile strength loss (measured after heated at 100° C. for 10 minutes) of 0% to 7.7%.\nAs described above, according to each of the Examples of these documents, the electro-deposited copper foils manufactured by these methods have low-profile deposit surfaces. The quality level thereof is excellent as a conventional low-profile electro-deposited copper foil and may be effective for forming a fine-pitch wiring. It is also disclosed that more excellent mechanical strength than that of a conventional electro-deposited copper foil can be obtained. It should be noted that the term “low-profile” used in a copper foil for a printed wiring board means that the height between peak and valley of a bonding side of a copper foil to be bonded to an insulating layer-forming material is low.\n[Patent Document 1] Japanese Patent Laid-Open No. 2004-263289\n[Patent Document 2] Japanese Patent Laid-Open No. 2004-339558"}
{"text": "1. Field of the Invention\nThe invention relates to a circuit breaker, and more particularly to a built-in intelligent circuit breaker capable of realizing automatic closing function, and a miniature intelligent circuit breaker with automatic closing function.\n2. Description of the Related Art\nAs to the electric power department, a circuit breaker is an essential device for ensuring the electricity safety and circuit switching.\nAs shown in FIG. 1A, an example of a circuit breaker in the prior art, which is called a circuit breaker body in the invention. The circuit breaker can be a three-phase or a single-phase circuit breaker, and can also be an ordinary miniature circuit breaker or a big switch. Each circuit breaker includes a box body; the box body includes an upper cover and a bottom box, wherein a notch groove is formed on the upper cover, and two ends of the notch groove are positioned in the on/off position correspondingly to the circuit breaker; a poke rod extends out of the notch groove; a circuit breaker actuating mechanism, a wire inlet end and a wire outlet end are arranged on the bottom box, and the circuit breaker actuating mechanism is triggered through the poke rod; the circuit breaker actuating mechanism further includes a linkage component connected with the poke rod, a movable contact arm is hinged at the lower end of a linkage rod of the linkage component, when the linkage component rotates under the action of the poke rod, the movable contact arm is driven to rotate, a moving contact on the movable contact arm is in contact with a stationary contact on a static contact piece, and the static contact piece is connected with the wire outlet end on the bottom box, so that the purpose of transmitting current out is achieved.\nAs shown in FIG. 1B and FIG. 1C, examples of a circuit breaker in the prior art, which are ordinary miniature circuit breakers. Each circuit breaker includes a box body; the common box body includes an upper cover and a bottom box, wherein a notch groove 1111 is formed on the upper cover, and two ends of the notch groove 1111 are positioned in the on/off position correspondingly to the circuit breaker; a handle 114 extends out of the notch groove; a circuit breaker actuating mechanism, a wire inlet end and a wire outlet end are arranged on the bottom box, and the circuit breaker actuating mechanism is triggered through the handle; the circuit breaker actuating mechanism further includes a linkage rod connected with the handle, a movable contact arm is hinged at the lower end of the linkage rod, a spring is arranged at the small end of the linkage rod, when the linkage rod rotates under the action of the handle, the movable contact arm is driven to rotate, a moving contact on the movable contact arm is in contact with a stationary contact on a static contact piece, and the static contact piece is connected with the wire outlet end on the bottom box, so that the purpose of transmitting current out is achieved; the movable contact arm is connected with a leading-in metal sheet through a metal connecting rod; an arc-extinguishing device is arranged at the lower part of the contact area between the moving contact and the stationary contact; meanwhile, a transverse partition board is arranged on the bottom box.\nHowever, the conventional circuit breaker still has the following defects:                1. In a prepayment system, the circuit breaker cannot be automatically closed after the electricity bills are paid by electricity consumers, and the electricity consumers have to deal with it by themselves, thus inconvenience is brought to the electricity consumers;        2. When the temperature is too high, and the use of electricity is abnormal, the opening speed of the conventional circuit breaker is slower; and        3. The circuit breaker can be closed by users under the condition without eliminating the condition of abnormal use of electricity, thus potential safety hazards may be caused.        "}
{"text": "A need exists for a method for offshore transfer of fluid that can form a hydrocarbon vapor, while maintaining a stable distance between a fixed or floating vessel and a transport ship safely without creating emissions to the environment.\nA need exists for a method that can be used to dynamically react to environmental conditions, such as wind and waves, to extend and retract a jib to maintain a stable distance between a fixed or floating vessel and the transport ship while offloading.\nA need exists for a method for maintaining a nominal distance between the fixed or floating vessel and the transport ship, and providing a quick release and emergency disconnect without spilling any of the fluid into the sea when the jib has extended to a maximum extension or retracted to a minimum retraction.\nA further need exists for a method that can safely transport personnel and equipment while floating in rough seas, while offloading a fluid that forms a hydrocarbon vapor, and while maintaining a transport ship at a nominal distance adjustably from the fixed or floating vessel.\nThe present embodiments meet these needs.\nThe present embodiments are detailed below with reference to the listed Figures."}
{"text": "Bus-masters are devices that can initiate a read or write transaction with the system memory of a bus-based computer system. An example of one such device is a processor. Processors retrieve instructions from the system memory and read, process, and write data to the system memory as dictated by the retrieved instructions.\nOne way to increase the performance of a computer system is to increase the rate data is transferred between the system memory and the bus-masters in the system. This can be accomplished by using high-speed memory chips for the system memory. However, the cost of memory increases dramatically with the speed of the memory. Thus, it is rarely practical to use the fastest memory available, especially in systems which require large quantities of system memory. Consequently, it has been discovered that by using a relatively small bank of relatively high-speed memory (\"cache memory\") as a buffer between a bus-master and the larger, slower bank of system memory, system performance can be greatly improved.\nThe data transfers between the system memory and the cache memory of a bus-master are controlled by a cache controller. A cache controller manages the cache memory based upon the needs of the bus-master which it serves. For example, if a processor requires data that is not currently stored in the processor's cache memory, the data must be transferred from the slower system memory to the cache memory of the processor. Once the data resides in the processor's cache memory, the processor may access the data for subsequent operations at the faster rate of the cache memory. Further, if the processor requires information that does not reside in cache memory and all of the cache memory blocks are already allocated, then the cache controller causes the cache memory to free up storage space for the newly requested data, either by writing over unmodified data in the cache memory or by storing modified data back to the system memory.\nIn systems having more than one cache, advantages may be gained by sharing one or more cache lines among multiple cache memories. However, such sharing could lead to data conflicts when multiple cache memories assert ownership over the same cache line at the same time. To avoid such conflicts, a protocol must be implemented to allocate the ownership of the cache lines among the various cache memories. Such protocols are referred to herein as ownership protocols, and generally allow the transfer of ownership of a cache line from one cache memory to another.\nOwnership protocols may be adapted to govern cache line ownership on systems having more than one cache-equipped bus-master, such as a system with a cache-equipped processor and a cache-equipped I/O device. However, in such a system, multiple versions of the same data may exist. For example, both the cache memory of the processor and the cache memory of the I/O device may hold copies of data originating from the same block in the system memory. Data that is accessible by more than one cache-equipped bus-master is referred to as \"shared data\". If either the I/O device or the processor modifies a block of shared data in its local cache memory, all other copies of the data, including any copy of the data in the cache of the other bus-master, become \"stale\" (no longer current).\nTo prevent the use of stale data, the bus-master that modifies a block of shared data must broadcast the modified version of the data to all other cache-equipped bus-masters in the system. However, this broadcasting process inherently requires a large amount of inter-device communication. Such communication is typically performed over the same bus which connects the bus-masters with the system memory. As a result, these busses can become crowded, resulting in slower data transmission rates and diminished system performance.\nIn view of the foregoing, a device and method for making a cache-equipped bus-master aware of whether any other cache-equipped bus-master is present in the system is clearly desirable. Further, a device and method which allows a cache-equipped bus-master to determine whether other cache-equipped bus-masters are on the system without increasing the data traffic on the bus connecting the bus-masters to the system memory is clearly desirable. Finally, a device for determining the shared status of a data block without increasing the data traffic on the bus connecting the bus-masters to the system memory is clearly desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to sports equipment and, particularly, to sports equipment stick handles.\n2. Description of the Background\nA sports equipment stick generally comprises an operational head joined to a handle on which the player places his hands. The head, for example, the head of a lacrosse stick or the blade of a hockey stick, refers generally to that part or parts of the sports equipment stick that is/are controlled by the handle. As used herein, “stick” refers to the stick as a whole, including the operational head and the handle.\nTo execute game skills, a player controls the sports equipment stick by grabbing and maneuvering the handle. In some sports, such as hockey and lacrosse, this control is referred to as “stick handling,” and game skills that involve stick handling include passing, shooting, checking, cradling, carrying, and scooping ground balls. For improved stick handling, it is generally desirable that the stick be light weight. A light weight stick provides improved ball sensitivity, allows quicker stick movement, reduces player fatigue, and provides better stick balance. The stick's weight can be decreased by decreasing the weight of the handle. To decrease the handle weight, either a lighter material must be used or the wall thickness must be decreased. Both options entail compromises.\nDecreasing the wall thickness also decreases the handle's strength, rigidity, and durability. Conventional stick handles comprise hollow metal tubes of aluminum, titanium, or other suitable strong lightweight alloys. Lacrosse stick handles generally have an octagonal cross section with uniform wall thickness. Given the octagonal shape, handle-on-handle impacts are almost always edge-to-edge impacts (it is rare that the impacting handles are perfectly aligned to result in a flat-side-surface impact). While executing game skills, it is common for one player's stick to strike another player's stick, whether intentionally or inadvertently. For example, a player executes a stick check to dislodge a ball from an opponent by intentionally striking the opponent's stick with the player's own stick. Often times, these strikes result in handle-on-handle impacts. The handle must have a substantial degree of strength, rigidity, and durability to withstand these impacts and resist dents (that lead to handle failure) while also minimizing stick weight for stick handling.\nThere is therefore a need in the art for an improved sports equipment stick handle that has the necessary strength to withstand handle-on-handle impacts, but also minimizes the handle's weight for superb stick handling. There is also a need in the art for an improved sports equipment stick handle that has increased strength to withstand handle-on-handle impacts while maintaining a similar weight to conventional handles."}
{"text": "The present invention relates to an improved trimmer head for use in flexible line rotary trimming devices used to trim grass, weeds and other vegetation. More particularly, the invention is directed to an improvement in the line loading of both “bump-feed” type trimmer heads such as those disclosed in U.S. Pat Nos. 4,458,419 and 4,959,904 and the more simple manually operated heads such as that disclosed in U.S. Pat. No. 4,145,809, the contents of said patents being incorporated herein by reference as though fully set forth below.\nTrimmer heads used in flexible line rotary trimmers generally carry one or two lengths of flexible nylon cutting line wrapped about an interior spool with the ends of the line or lines projecting outwardly through opposed apertures in the side wall of the trimmer head. The head is threadably mounted on the end of an elongated shaft and rotated at a high velocity by a gas or electric motor so that the ends of the cutting line project radially from the head and sever weeds or other vegetation. When cutting line projecting from the head breaks off or becomes overly worn, it must be severed and fresh line extended from the spool through the line outlet eyelets in the side of the housing. Bump-feed type heads include a line feed-out mechanism which responds to a bump on the ground intentionally applied by the operator to feed out a measured length of fresh cutting line which is typically cut to the desired length by a knife blade projecting from a shield attached to the trimmer above the cutting head and spaced a predetermined distance from the perimeter of the trimmer head housing. Manual heads do not include any such line feed-out mechanism. A fastening nut that holds the housing portion of the trimmer head to the spool must be loosened so that the spool can be separated from the housing and manually rotated relative to the housing to pay out additional cutting line. The spool and housing are then re-secured by the fastening member.\nIn both bump-feed and manual heads, the length or lengths of cutting line are typically wound onto the spool by hand. As most cutting heads employ two lengths of line projecting from opposed sides of the cutting head, care must be taken during the winding of the spool to avoid crossing or otherwise tangling of the two lines within the spool which interferes with the paying out of fresh line. This is particularly important in bump-feed heads where centrifugal force is utilized to pull the new lengths of line from the spool during use as the head is being bumped against the ground as any line tangle will interfere with the proper feeding of the line. Difficulty in properly loading the line on the spool is the most common complaint of home users of flexible line trimmers. It also is a time consuming task for the professional user.\nThe early bump-feed mechanisms typically consisted of a dog or friction clutch located between the spool of line and the surrounding housing. By bumping an extension of the spool on the ground, or other fixed object, the friction clutch was temporarily disengaged for a length of time dependant on the duration of the bump. The dog clutch released by the bump then abruptly engaged at the next opportunity to feed out line in segment lengths which were related to the engagement points of the dog clutch. Such dog clutches had outwardly extending ribs which engaged inwardly extending abutment tangs and therefore depended upon a skillful bump when it was desired to deliver only one segment length. However, friction within such devices and overzealous bumping often resulted in two or more line segments being fed out, particularly especially when the device has been in use and the corners on the ribs and tangs became worn such that positive engagement was no longer assured. The unavoidable abrupt operation of the dog clutch caused such wear to occur.\nA bump-feed-out mechanism was subsequently developed that automatically fed out a predetermined length of line with each bump, regardless of the duration of the bump, and which did not lose this ability with extended use. That device is disclosed in and is the subject of the incorporated reference, U.S. Pat. No. 4,458,419. As described therein in detail, the improved trimmer head contained a spool holding one or more coils of cutting line and a simplified mechanism that selectively allowed relative movement of the spool with respect to the housing in response to bumping of the head on the ground to pay out measured lengths of line. The simplified pay-out mechanism included a novel spring-loaded cam and cam follower arrangement in which the cam follower included two pair of diametrically opposed and generally inwardly facing abutment surfaces arrayed about the axis of rotation of the trimmer housing. The abutment surfaces were thus spaced 90° apart and were carried by a depending cylindrical wall that circumscribed an interior chamber. The cam member was disposed within the chamber in threaded engagement with the extended lower end of the drive bolt of the trimmer housing and defined two vertically adjacent cams, each cam being of a square configuration and defining four perpendicularly disposed cam surfaces adapted to engage the abutment surfaces on the cam follower. The upper cam was rotationally offset 45° from the lower cam.\nIn operation, the housing was rotationally driven by the drive bolt through a connection between the upper end of the bolt and the trimmer drive means. The housing and cam member was thus driven by the drive bolt, which in turn drove the cam follower and the spool mounted thereon due to the engagement between the cam surfaces on the cam member and the abutment surfaces on the cam follower. The line carrying spool was disposed about the cylindrical wall of the cam follower and attached thereto via a pair of opposed outwardly projecting studs on the cam follower member that extend into slots formed in the inner portion of the spool. The spool was provided with a bumper at its lower end such that when the bumper was pressed against or bumped on the ground, the housing moved downwardly with respect to the spool against the force of a spring, disengaged the lower cam from the abutment surfaces on the cam follower and allowed the cam member to rotate 45° relative to the cam follower, whereupon the cam surfaces of the upper cam would abut the abutment surfaces on the cam follower. That imparted a similar degree of relative rotation between the spool and the housing. Once the force of the bump was dissipated, the spring loading forced the spool and housing back to their relative positions, which released the cam surfaces on the upper cam from the cam follower abutment surfaces and allowed another 45° of relative rotation of the cam member and cam follower and thus of the spool and housing, for a total of 90° of rotation per bump, which provided the predetermined relative rotation between the housing and spool needed to pay out a desired length of line through the apertures in the trimmer housing. Since the cams interacted with simple, inwardly facing cam follower surfaces formed only on a single level, the release mechanism was deemed relatively economical to manufacture and, due to the large abutment surface areas presented between the cams and cam follower, the device was durable, trouble free and reliable.\nBecause of early difficulties in molding some of the components of the cutting head disclosed in U.S. Pat. No. 4,458,419, the head became more expensive to manufacture than anticipated. New material developments subsequently reduced the cost of manufacture. In the meantime, however, a similar bump-feed drive mechanism was developed in which large square cams were formed on the upper and lower outer radial surfaces of the spool and the corresponding cam followers were formed by upper and lower portions of the housing which surrounded the spool. Such a head is disclosed in the incorporated reference, U.S. Pat. No. 4,959,904, and is still in production.\nOver the years, with increasing competition from offshore manufacturers, it became clear that even with the development of new materials the earlier bump-feed mechanism covered by U.S. Pat. No. 4,459,419 was not as economical to manufacture as earlier believed. It contained several parts, some of which had to be hand assembled. In addition, vibration, the threaded engagement between the cam member and the drive bolt, and the heat generated by the trimmer required the use of a chemical bonding agent having a high melting point to prevent the cam member and cam follower from breaking loose from the drive bolt. Such agents, however, had extremely high break way torques, rendering the threaded connection effectively permanent. As a result, certain components of the head could not be replaced when worn. Thus, that head was significantly modified so as to retain all of the advantages of its predecessor yet utilize fewer component parts and obviate the need for any hand assembly and use of chemical bonding. That modification is the subject of a pending U.S. patent application entitled “Trimmer Head for Use in Flexible Line Rotary Trimmers”, filed Oct. 2, 2003 and identified by Ser. No. 10/677,700, of which this application is a continuation-in-part. The bump-feed and manual heads of the present invention retain all of the advantages of the above-described heads and adds thereto the ability to far more quickly and easily uniformly wind lengths of cutting line onto the spool without materially increasing the cost of production."}
{"text": "In the world of telecommunications, the interaction of a plurality of telecom networks operated by more than one entity has become more prevalent. Calls, texts, and other types of communications may be between different users in many countries, almost anywhere in the world, and provide both voice and data access for such users.\nWhen a user initiates a communication or service that extends beyond his home network (i.e., the network to which the user subscribes) to an additional network, that additional network provides services by connecting the communication (e.g., “bridging”, routing, etc. the call). Such services of the additional network benefit both the users involved with the call/communication and the home network of the user. The home network benefits because its customers are able to engage in cross-network communications, thereby making those customers satisfied with the functionality of their home network. Typically, those customers receive one billing statement from their home network; the billing statement lists those services provided by the home network and beyond (e.g., long distance calling, additional services provided by additional network(s), value added services, etc.)\nIn order to provide the aforementioned billing statements, the various additional networks (if used) must be capable of working with the home network and sending the home network billing information (e.g., billing rates for the additional network, time increments for charges, access fees, etc.). One manner in which this is typically accomplished is that the home network periodically communicates with additional network(s) to obtain up-to-date billing information for the respective network. Additionally, any network(s) providing services for the home network may periodically send a report of the provided services and related costs to the home network. After receipt thereof, the home network may apply the charges to the applicable billing statement(s).\nUnfortunately, such periodic updating may involve various inefficiencies, such as, but not limited to, untimely receipt of billing amendments, new contract rate(s), relevant charges in the current billing cycle; delaying charges on a billing statement or invoice; sending more than one invoice to reflect charges from different networks for the same call; etc. Such inefficiencies lead to additional costs of administering such networks. For example, if an additional network connects a call from England to New York City, the additional network may charge ten dollars U.S. total. However, the additional network may not report the ten dollar charge for the call to the home network until after the home network sends a billing statement listing that call and the related charges to the customer. As such, a customer may be charged additionally for a phone call after already paying a previous billing statement or invoice for that same call (i.e., the customer thinks they paid for that service in full when they may have not). Alternatively, the home network may try to avoid confusion from the aforementioned scenario by delaying any charges to the customer for the call until it receives some type of billing information from the additional network(s) involved with the call. As a result, the home network may not receive payment from the customer for its provided services within the normally contracted timeline, thereby leading to the aforementioned inefficiencies. As such, there is a need in the art to prevent imperfect communication between networks, additional administration costs, etc.\nIn view of the above, it would be desirable to provide a system, method, and computer-readable storage medium that permits the outsourcing of telecom billing and other administrative services, so that a network servicing or interfacing with a second network can provide billing and other administrative services to said second network. It would also be desirable for a telecom wholesale provider to be able to provide billing and other administrative services to another telecom provider to which the wholesale provider provides services. Additionally, whether or not any network provides telecom services to another network, outsourcing by one telecom provider of its administrative and billing functions to another entity is often desirable for business reasons."}
{"text": "S-Adenosylmethionine (SAMe) is found in almost every tissue and fluid in the body. SAM plays a crucial role in the process called transmethylation. Methylation is involved in nearly every aspect of life. SAM is the primary “methyl” donor for a variety of methyl-transfer reactions in DNA, RNA, proteins, lipids, and small molecules in the body. Proper DNA methylation is essential for normal embryonic development. Methyl-transferase gene homozygously deleted (knocked out) has been proven lethal (Pegg, A. E., Feith, D. J., Fong, L. Y., Coleman, C. S., O'Brian, T. G., and Shantz, L. M., 2003, Biochem. Soc. Trans. 31, 356-360). DNA improperly methylated has been found in many tumors. Alterations in DNA methylation patterns induce the expression of oncogens or silence the expression of tumor suppressor genes, and methyl deficient diets have been shown to promote liver cancer in rodents.\nThe transsulfuration begins with S-adenosylhomocysteine (SAH), the residual structure of SAM upon donating the methyl group (transmethylation). Hydrolysis of SAH yields homocysteine, which in turns converts to cystathionine, then cysteine, and eventually, to glutathione, the hepatocellular antioxidant and life-saving detoxification agent.\nThe aminopropylation is another process initiated with SAM through decarboxylation. The decarboxylated SAM then couples with putrescine to generate spermidine and spermine which are critical to cell growth, differentiation and the stability of DNA and RNA. Furthermore, Methylthioadenosine (MTA), the by-product of polyamine synthesis, is a powerful analgesic and anti-inflammatory agent. This may be, at least partially, responsible for the clinical benefits observed in the treatment of osteoarthritis, rheumatoid arthritis and fibromyalgia with SAMe.\nSAMe plays a role in the immune system, maintains cell membranes, and helps produce and break down brain chemicals, such as serotonin, melatonin, and dopamine. Deficiency of either vitamin B12 or foliate can reduce the level of SAMe. SAMe is also an antioxidant, a substance that protects the body from damaging reactive oxygen molecules in the body. These reactive oxygen molecules can come from inside the body or from environmental pollution and are thought to play a role in the aging process and the development of degenerative disease. In general, SAMe is thought to raise the level of functioning of other amino acids in the body.\nBy way of further background, S-adenosyl-1-methionine is a substrate of an enzyme lyase that converts S-adenosyl-1-methionine to the molecule methylthioadenosine and homoserine; it is an aminobutyric chain donor to tRNA; it is an aminoacidic chain donor in the biosynthesis of biotin; SAM-e, after decarboxylation, is the donor of aminopropyl groups for the biosynthesis of neuroregulatory polyamines spermidine and spermine. (Zappia et al (1979), Biomedical and Pharmacologcial roles of Adenosylmethionine and the Central Nervous System, page 1, Pergamon Press. N.Y.)\nSAM-e has been used clinically in the treatment of liver disease (Friedel H, Goa, K. L., and Benfield P., (1989), S-Adenosyl-1-methionine: a review of its pharmacological properties and therapeutic potential in liver dysfunction and affective disorders in relation to its physiological role in cell metabolism. Drugs. 38, 389-416), arthritis (Di Padova C, (1987), S-adenosyl-1-methionine in the treatment of osteoarthritis: review of the clinical studies. Am J. Med. 83, (Suppl. 5), 6-65), and depression (Kagan, B, Sultzer D. L., Rosenlicht N and Gerner R. (1990), Oral S-adenosylmethionine in depression: a randomized, double blind, placebo-controlled trial. Am. J. Psychiatry 147, 591-595.) Alzheimer's patients have reduced cerebral spinal fluid levels of S-adenosyl-1-methionine (Bottiglieri et al, (1990), Cerebrospinal fluid S-adenosyl-1-methionine in depression and dementia: effects of treatment with parenteral and oral 5-adenosyl-1-methionine. J. Neurol. Neurosurg. Psychiatry 53, 1096-1098.) In a preliminary study, SAM-e was able to produce cognitive improvement in patients with Alzheimer's disease. (Bottiglieri et al (1994), The clinical potential of admetionine (S-adenosyl-1-methioinine) in neurological disorders. Drugs 48, 137-152.) SAM-e brain levels in patients with Alzheimer's disease are also severely decreased. (Morrison et al, (1996), Brain S-adenosylmethionine levels are severely decreased in Alzheimer's disease, Journal of Neurochemistry, 67, 1328-1331.) Patients with Parkinson's disease have also been shown to have significantly decreased blood levels of SAM-e. (Cheng et al, (1997), Levels of L-methionine S-adenosyltransferase activity in erythrocytes and concentrations of S-adenosylmethionine and S-adenosylhomocysteine in whole blood of patients with Parkinson's disease. Experimental Neurology 145, 580-585.)\nSAM-e levels in patients treated with the antineoplastic drug methotrexate are reduced. Neurotoxicity associated with this drug may be attenuated by co-administration of SAM-e. (Bottiglieri et al (1994), The Clinical Potential of Ademetionine (S-adenosylmethionine) in neurological disorders, Drugs, 48 (2), 137-152.)\nCerebral spinal fluid levels of SAM-e have been investigated in HIV AIDS dementia Complex/HIV encephalopathy and found to be significantly lower than in non-HIV infected patients. (Keating et al (1991), Evidence of brain methyltransferase inhibition and early brain involvement in HIV positive patients Lancet: 337:935-9.)\nDe La Cruz et al have shown that SAM-e, chronically administered, can modify the oxidative status in the brain by enhancing anti-oxidative defenses. (De La Cruz et al, (2000), Effects of chronic administration of S-adenosyl-1-methionine on brain oxidative stress in rats. Naunyn-Schmiedeberg's Archives Pharmacol 361: 47-52.) This is similar to results obtained with SAM-e in liver and kidney tissue. Thus SAM-e would be useful as an antioxidant.\nOral SAM-e administration to patients with and without liver disease has resulted in increases in liver glutathione levels. (Vendemiale G et al, (1989), Effect of oral S-adenosyl-1-methionine on hepatic glutathione in patients with liver disease. Scand J Gastroenterol; 24: 407-15. Oral administration of SAM-e to patients suffering from intrahepatic cholestasis had improvements in both the pruritus as well as the biochemical markers of cholestasis. (Giudici et al, The use of admethionine (SAM-e) in the treatment of cholestatic liver disorders. Meta-analysis of clinical trials. In: Mato et al editors. Methionine Metabolism: Molecular Mechanism and Clinical Implications. Madrid: CSIC Press; 1992 pp 67-79.) Oral SAM-e administration to patients suffering from primary fibromyalgia resulted in significant improvement after a short term trial. (Tavoni et al, Evaluation of S-adenosylmethioine in Primary Fibromaylgia. The American Journal of Medicine, Vol 83 (suppl 5A), pp 107-110, 1987.) SAM-e has been used for the treatment of osteoarthritis as well. (Koenig B. A long-term (two years) clinical trial with S-adenosylmethionine for the treatment of osteoarthritis. The American Journal of Medicine, Vol 83 (suppl 5A), Nov. 20, 1987 pp 89-94)\nSAM-e is clinically useful in many apparently unrelated areas because of its important function in basic metabolic processes. One of its most striking clinical uses is in the treatment of alcoholic liver cirrhosis that, until now, remained medically untreatable. Mato et al demonstrated the ability of oral SAM-e in alcoholic liver cirrhosis to decrease the overall mortality and/or progression to liver transplant by 29% vs 12% as compared with a placebo treated group. (Mato et al (1999), S-adenosylmethionine in alcohol liver cirrhosis: a randomized, placebo-controlled, double blind, multi-center clinical trial, Journal of Hepatology, 30, 1081-1089.)\nSam-e also attenuates the damage caused by tumor necrosis factor alpha and can also decrease the amount of tumor necrosis factor alpha secreted by cells. Consequently, conditions in which this particular inflammatory factor is elevated would benefit from the administration of SAM-e. (Watson W H, Zhao Y, Chawla R K, (1999) Biochem J August 15; 342 (Pt 1):21-5. S-adenosylmethionine attenuates the lipopolysaccharide-induced expression of the gene for tumour necrosis factor alpha.) SAM-e has also been studied for its ability to reduce the toxicity associated with administration of cyclosporine A, a powerful immunosuppressor. (Galan A, et al, Cyclosporine A toxicity and effect of the s-adenosylmethionine, Ars Pharmaceutica, 40:3; 151-163, 1999.)\nSAM-e, incubated in vitro with human erythrocytes, penetrates the cell membrane and increases ATP within the cell thus restoring the cell shape. (Friedel et al, S-adenosyl-1-methionine: A review of its pharmacological properties and therapeutic potential in liver dysfunction and affective disorders in relation to its physiological role in cell metabolism, Drugs 38 (3):389-416, 1989)\nSAM-e has been studied in patients suffering from migraines and found to be of benefit. (Friedel et al, S-adenosyl-1-methionine: A review of its pharmacological properties and therapeutic potential in liver dysfunction and affective disorders in relation to its physiological role in cell metabolism, Drugs 38 (3): 389-416, 1989)\nSAM-e has been administered to patients with peripheral occlusive arterial disease and was shown to reduce blood viscosity, chiefly via its effect on erythrocyte deformability.\nSAM-e is commercially available using fermentation technologies that result in SAM-e formulations varying between 60 and 80% purity. (That is, the final product contains 60-80% of the active or (S, S)-SAM-e and 20-40% of the inactive or (R, S)-SAM-e.) (Gross, A., Geresh, S., and Whitesides, Gm (1983) Appl. Biochem. Biotech. 8, 415.) Enzymatic synthetic methodologies have been reported to yield the inactive isomer in concentrations exceeding 60%. (Matos, J R, Rauschel F M, Wong, C H. S-Adenosylmethionine: Studies on Chemical and Enzymatic Synthesis. Biotechnology and Applied Biochemistry 9, 39-52 (1987). Enantiomeric separation technologies have been reported to resolve the pure active enantiomer of SAM-e. (Matos, J R, Rauschel F M, Wong, C H. S-Adenosylmethionine: Studies on Chemical and Enzymatic Synthesis. Biotechnology and Applied Biochemistry 9, 39-52 (1987; Hoffman, Chromatographic Analysis of the Chiral and Covalent Instability of S-adenosyl-1-methionine, Biochemistry 1986, 25 4444-4449: Segal D and Eichler D, The Specificity of Interaction between S-adenosyl-1-methionine and a nucleolar 2-0-methyltransferase, Archives of Biochemistry and Biophysics, Vol. 275, No. 2, December, pp. 334-343, 1989) Newer separation technologies exist to resolve enantiomers on a large commercial production scale at a very economic cost. In addition, it would be conceivable to synthesize the biologically active enantiomer using special sterioselective methodologies but this has not been accomplished to date.\nDe la Haba first showed that the sulfur is chiral and that only one of the two possible configurations was synthesized and used biologically. (De la Haba et al J. Am. Chem. Soc. 81, 3975-3980, 1959) Methylation of RNA and DNA is essential for normal cellular growth. This methylation is carried out using SAM-e as the sole or major methyl donor with the reaction being carried out by a methyltransferase enzyme. Segal and Eichler showed that the enzyme bound (S, S)-SAM-e 10 fold more tightly than the biologically inactive (R, S)-SAM-e thus demonstrating a novel binding stereospecificity at the sulfur chiral center. Other methyltransferases have been reported to bind (R, S)-SAM-e to the same extent as (S, S)-SAM-e and thus (R, S)-SAM-e could act as a competitive inhibitor of that enzyme. (Segal D and Eichler D, The Specificity of Interaction between S-adenosyl-1-methionine and a nucleolar 2-0-methyltransferase, Archives of Biochemistry and Biophysics, Vol. 275, No. 2, December pp. 334-343, 1989; Borchardt R T and Wu Y S, Potential inhibitors of S-adenosylmethionine-dependent methyltransferases. Role of the Asymmetric Sulfonium Pole in the Enzymatic binding of S-adenosyl-1-methionine, Journal of Medicinal Chemistry, 1976, Vol 19, No. 9, 1099-1103.)\nSAM-e (whether in its optically pure enantiomeric form or in an enantiomeric or racemic mixture) presents certain difficult problems in terms of its stability at ambient temperature that result in degradation of the molecule to undesirable degradation products. SAM-e (and thus its enantiomers) must be further stabilized since it exhibits intramolecular instability that causes the destabilization and breakdown of the molecule at both high as well as ambient temperatures. SAM-e has therefore been the subject of many patents directed both towards obtaining new stable salts, and towards the provision of preparation processes that can be implemented on an industrial scale. The present patent thus envisions the use of any of the salts of SAM-e already disclosed in the prior art to stabilize the enantiomeric forms of SAM-e.\nThe clinical diagnostic field has seen a broad expansion in recent years, both as to the variety of materials of interest that may be readily and accurately determined, as well as the methods for the determination. Over the last several decades, testing for numerous substances such as drugs of abuse, or other biological molecules of interest has become commonplace. In recent years, immunoassay based on the interaction of an antibody with an antigen has been extensively investigated for this purpose. Based on the unique specificity and high affinity of antibodies, an immunoassay can accurately and precisely quantitate substances at the very low concentrations found in biological fluids.\nAccordingly, there is a need for improved methods of detection and diagnosis of cancer and other diseases, as well as methods for monitoring the progress of the diseases and monitoring the progress of various treatments for cancer and other diseases by quantitating the methylating index as a biomarker."}
{"text": "Sales of cylinder positioning systems have generally fallen into one of two categories, simple and complex. A simple system is inexpensive, and typically uses timing or limit switches to achieve control. Although this type of system has significant cost advantages, it has some performance disadvantages. Limit switch systems lack the ability to dynamically change the point or points at which the cylinder stops, and timing controlled systems require constant pressure, load, and wear. Complex systems are more tolerant of pressure, load, and wear changes but have disadvantages in cost and complexity. Complex systems can cost between 10 and 20 times more than the simple system. The complex systems typically use fragile or expensive sensor technology, and require a trained, experienced person for installation."}
{"text": "1. Field of the Invention\nThe present invention relates to a dielectric ceramic and electronic parts using the same, and particularly, to the dielectric ceramics which are applied, for example, in a micro-stripline device, a monolithic IC, a polarizer and so on, and the electronic parts using the dielectric ceramics.\n2. Description of the Prior Art\nConventionally, as a dielectric ceramic used in electronic parts, those having a uniform composition and a uniform dielectric constant distribution have been used.\nFor example, as a dielectric substrate used in a conventional micro-stripline device, there was the dielectric substrate having a uniform composition structure in which such material as alumina and fluororesin having a dielectric constant of 10 or so are used.\nIn such a dielectric substrate, by increasing the dielectric constant, the micro-stripline device or the substrate can be miniaturized. However, when the dielectric constant is increased, a width of stripline must also be narrowed in order to satisfy a characteristic impedance standard.\nWhen considering miniaturization of the stripline device with the conventional dielectric substrate, under the condition of 1 mm substrate thickness and 50 .OMEGA. characteristic impedance, a necessary line width of about 1.0 mm for the substrate of dielectric constant of 10. On the other hand, a necessary line width for the substrate of dielectric constant of 30 is about 0.2 mm. Accordingly, when miniaturizing by increasing the dielectric constant of the substrate, the line width becomes so narrow that it is difficult to contact to external circuits at input/output terminals.\nSimilarly, when the substrate having the uniform composition is used as the monolithic IC substrate, the line width formed on the substrate is determined by the characteristic impedance. Therefore, when the dielectric constant of the substrate is increased to reduce the line width, it is difficult to make contact to the external circuits.\nWhile, in a satellite broadcasting system, as shown in FIGS. 25(A) and (B), a polarizer comprising a waveguide 1 into which a dielectric plate 2 having a dielectric constant of about 10 is inserted was used. This dielectric plate 2 is that, the dielectric constant of the dielectric plate in an axial direction of the waveguide is changed continuously by forming V-shaped inclined portions 3 at its opposite ends. As such, by changing the dielectric constant of the dielectric plate, when an electric wave is incided to the polarizer or radiated from the polarizer, reflection of the electric wave at a border of the dielectric plate and air is prevented.\nAlso, in this polarizer, it can be miniaturized by using the ceramics having the high dielectric constant for the dielectric plate. However, when the ceramics having the high dielectric constant is used, a difference of dielectric constant between air and the dielectric plate becomes larger, and a propagation wave such as an interface wave is produced at the border between the air and dielectric plate, results in increasing a transmission loss."}
{"text": "The present invention relates to the structure of a casing of a miniature portable apparatus and a clip attachable thereto and, more particularly, to a structure which allows a clip to be firmly affixed to the casing of a radio pager while promoting easy attachment and detachment of the clip.\nA wide variety of miniature portable apparatuses are extensively use today and include radio pagers and audio apparatus. A radio pager, for example, may be provided with a clip to be put in the user's chest pocket or on the user's waist belt. The clip is attached to the casing of such a radio pager. Specifically, the casing is usually made of plastic and provided with an elongate groove on the top thereof and an elongate lug on the bottom, while an elastic bate plate is made of metal and provided with an engaging portion engageable with the groove in a upper end portion thereof and a slot engageable with the lug in a generally J-shaped lower end portion. The clip has a clip body which is rotatably supported by the base plate. To attach the clip to the casing, the engaging portion of the base plate is put in the groove of the casing, and then the base plate is rotated about the groove. As a result, the slot of the base plate mates with the lug of the casing due to the resiliency of the base plate. To remove the clip from the casing, the lower end portion of the base plate is urged away from the casing by, for example, a tool to release the slot from the lug.\nHowever, the problem with the conventional clip structure is that the edge of the lower end of base plate moves in contact with the casing in the event when the clip is attached to or detached from the casing, scratching or even breaking the casing. Another problem is that the resiliency of the base plate used to attach and detach the clip from the casing is apt to differ from one base plate to another for production reasons. Therefore, some base plate come off the casing rather easy, while other are difficult to remove from the casing. In addition, since the edge of the base plate is forcibly urged away from the casing by a tool, it is likely that the casing is scratched or the base plate is deformed."}
{"text": "Early aircraft AUTOBRAKE systems controlled airplane deceleration to one of several predetermined aircraft deceleration settings. Thus, if a pilot wished to stop the airplane at a certain point on the runway, for example at a runway exit, it is unlikely that the AUTOBRAKE system would provide a deceleration setting that matched the stopping distance to the selected runway stop point. The pilot's choice, then, was to select one of the deceleration settings, and, if the airplane was decelerating too quickly, the pilot can select another AUTOBRAKE setting or disarm the AUTOBRAKE system and use pedal braking to avoid stopping short of the desired point. If, however, the airplane was decelerating too slowly, the pilot would again need to revert to a different AUTOBRAKE setting or to pedal braking to increase deceleration to stop the airplane at the selected stopping point. Both of the above conditions resulted in uneven deceleration that is apparent to the airplane passengers.\nAircraft AUTOBRAKING systems have been developed which apply predetermined braking to an aircraft. These AUTOBRAKING systems typically provide four or five selectable deceleration rates which can be used to decelerate the landing aircraft at the deceleration rate selected by the pilot. The result being either the braking of the aircraft at a specific rate until stopped or the pilot taking command of braking before the aircraft stops.\nModern aircraft can have the ability to determine position on the runway, calculate the deceleration rate for a specific exit taxiway on that runway, and automatically modulate the application of the brakes so as to arrive at the exit taxiway at a speed slow enough as to exit at the predefined exit taxiway."}
{"text": "Field of the Invention\nThe present invention relates to a communication apparatus provided with a plurality of network interfaces, a control method therefor, and a storage medium storing a control program therefor.\nDescription of the Related Art\nIn recent years, a communication apparatus that can search for other devices on a network is developed in order to provide a service while devices such as image forming apparatuses cooperate mutually. The communication apparatus generally uses a method using a multicast packet or a broadcast packet (for example, see Japanese Laid-Open Patent Publication (Kokai) No. 2000-181849 (JP 2000-181849A)).\nThe multicast packet is an IP (Internet Protocol) packet used when transmitting the same data to specific devices that stand on a network. One of the features of the multicast packet is that a TTL (TimeToLive) value is included in a packet. The TTL value expresses life time of the IP packet, and the value decreases one by one when passing through a router. That is, when the TTL value is set to “1” or more at a transmitting source of the multicast packet, it becomes possible to transmit the same multicast packet to networks (subnets) in a wide area beyond a router.\nOn the other hand, the broadcast packet is an IP packet used when transmitting the same data to many and unspecified devices that stand on a network. The broadcast packet cannot exceed a router, but can be transmitted only to devices on the network (subnet) that is formed by the router. As compared with this, a packet that specifies a single device that stands on a network and is transmitted to an IP address of the device is called a unicast packet.\nConventionally, a multi-home device that has a plurality of network interfaces (I/F) sends out a multicast packet or a broadcast packet from the network interfaces when searching for other devices on a network. In this case, the multi-home device of which the network interfaces belong to the same subnet has a problem in that device-search responses are received in duplication from devices that stand on the subnet.\nSince the multi-home device of which the network interfaces connect with different subnets becomes possible to search beyond a subnet when searching with the multicast packet of which the TTL value is set to “1” or more, the same problem occurs."}
{"text": "1. Field of the Invention\nThe present invention relates to a moving mechanism, for example, a moving mechanism that moves a form measurer.\n2. Description of Related Art\nAn apparatus shown in FIG. 13 is known, for example, as a form measuring apparatus 20 measuring a profile shape of a measured object W (Japanese Patent No. 2539998). A support pillar 2 stands on a base 1. A moving screw is rotatably provided inside the support pillar 2. An operation handle 2A is provided at an upper end of the moving screw (the moving screw is hidden inside the support pillar and thus invisible in FIG. 13). A lift member 3 is provided to the support pillar 2 so as to be movable upward and downward. A nut fixated to the lift member 3 is threaded together with the moving screw. Thus, rotation of the moving screw raises and lowers the nut and the lift member 3 (the nut is hidden inside the lift member 3 and thus invisible in FIG. 13).\nA form measurer 11 is provided to the lift member 3 via a swing member 4. The swing member 4 is angle-adjustable with reference to a swing axis, which is an axis orthogonal to an axis of the support pillar 2. The form measurer 11 swingably supports a measuring arm 12 and electrically detects a swing amount of the measuring arm 12. A stylus 13 projects from a fore end of the measuring arm 12, the stylus 13 moving along a surface profile of the measured object W.\nFor measurement, the measured object W is placed on the base 1 and the stylus 13 is positioned at a measurement start point of the measured object W. At this time, the form measurer 11 needs to be adjusted to be positioned such that the stylus 13 is positioned slightly higher than a height of a measured surface of the measured object W. Then, a user rotates the operation handle 2A to raise and lower the lift member 3 by way of screw movement.\nSince fine positioning is required for the form measuring apparatus 20, a screw lead of the moving screw is designed to be fine. This is beneficial to enable fine positioning; however, a circumstance concurrently arises where a movement amount per rotation of the operation handle 2A is extremely limited. Given that the form measurer 11 needs to be moved substantially every time the measured object W is replaced, positioning of the form measurer 11 requires a substantial amount of time.\nIn view of the above, the inventors of the present application have already devised a moving mechanism using twist rollers, instead of screw movement (Japanese Patent Nos. 5049514 and 5235039).\nThe moving mechanism disclosed in Japanese Patent Nos. 5049514 and 5235039 has a plurality of twist rollers capable of making contact with and separating from a drive shaft. A collet holding and releasing the drive shaft is provided to pivotally support the twist rollers. When the collet holds the drive shaft, the twist rollers come into contact with the drive shaft. When the collet releases the drive shaft, the twist rollers separate from the drive shaft. Expansion and contraction of the collet is controlled by a cam. Specifically, a cam having an elliptical shape in cross section is inserted to the collet. Rotating the cam changes a gap of the collet. Rotating the drive shaft in a state where the twist rollers are in contact with the drive shaft allows advancement and retraction by an amount of a twist angle of the twist rollers. In other words, fine movement is achieved. Meanwhile, in a state where the twist rollers are separated from the drive shaft, the collet is free and thus high-speed movement (rough movement) is enabled. With a moving body, such as a form measurer, provided to the collet, for example, switching between fine movement and rough movement of the form measurer (moving body) is enabled.\nThe configurations of Japanese Patent Nos. 5049514 and 5235039 reliably allow high-speed movement (rough movement) and fine positioning (fine movement). However, the configurations slightly lack in operability. In the configurations of Japanese Patent Nos. 5049514 and 5235039, for example, the cam needs to be rotated every time the fine movement and the rough movement are switched. When the fine movement and the rough movement are switched repeatedly to adjust positioning of the form measurer (moving body), it is cumbersome to operate the cam every time.\nIn addition, a practical circumstance as below arises. For example, the cam is rotated to separate the twist rollers from the drive shaft. At this time, the collet is free and thus the high-speed movement (rough movement) is enabled. However, when the user's hand is removed, the collet may operate in a manner not intended by the user. This circumstance is not notable in horizontal movement as described in Japanese Patent Nos. 5049514 and 5235039. In a vertically moving mechanism, however, in which a drive shaft is provided vertically, a collet may drop together with a form measurer (moving body). Thus, the user needs to support the form measurer (moving body) with one hand and then rotate the cam with the other hand to bring the twist rollers into contact with the drive shaft. Such an operation is very cumbersome and requires the user to perform in an unnatural manner.\nSuch a circumstance is not limited to the form measurer and is common to a moving mechanism that requires fine positioning and is sometimes required to make a large movement as well."}
{"text": "The invention relates to a burner for a rotary tubular furnace for the production of cement clinker, with a duct, arranged annularly within a burner carrier tube, for the discharge of primary fuels, for example coal dust, with at least one annular primary-air duct and with at least one tube, integrated in the burner, for the pneumatic transport of particulate secondary fuels and for the blowing these out at the burner mouth.\nIn a cement clinker production line, calcined raw cement meal is burnt in the sintering zone of a rotary tubular furnace into cement clinker. To heat the rotary tubular furnace, a long burner lance is introduced into the furnace outflow end through the stationary furnace outflow housing, at the mouth of which burner lance the fuels introduced into the lance burn so as to form a burner flame. The correct temperature, the length and the other configuration of the burner flame are important in the formation of clinker minerals in the rotary tubular furnace. Development tends to give the rotary tubular furnace itself as short a build as possible by virtue of high-quality calcination of the raw cement meal outside the rotary tubular furnace, so that, in reaction to this, the burner flame is, as a rule, to be as short and as hot as possible. Increasingly often, instead of liquid and gaseous fuels, the fuels used are solid fuels, in particular coal dust, but, recently, also pneumatically transportable particulate waste fuels, such as, for example, waste plastic granulates, etc., as secondary fuels.\nKnown rotary-furnace burners are often designed as what are known as three-duct burners (for example DE 43 19 363 A1), with at least three ducts concentric to one another, that is to say the pneumatically transported coal dust flows as fuel through the middle burner duct and emerges through an annular gap nozzle, the outflowing coal dust being surrounded by radially inner and by radially outer primary air as combustion air. The radially outer air, also called jet air, is subdivided, by means of a multiplicity of individual nozzles arranged in the annular jet-air duct, into a large number of individual high-velocity primary-air jets which generate a vacuum zone in their surroundings, that is to say the many high-velocity primary-air jets serve as propulsive jets on the injector principle, by virtue of which the large mass of the virtually stationary hot secondary air of, for example, about 1000° C., which surrounds the rotary-furnace burner, is sucked inwards in the direction of the core of the burner flame, where an intensive intermixing of the hot secondary air with the coal dust emerging through the annular gap nozzle takes place, the intention being that the coal dust should bum quickly and completely so as to form a short hot flame.\nIn the known three-duct burner, it would not be possible to inject particulate secondary fuel through the coal dust duct by means of its annular gap nozzle, since this secondary fuel would block the annular gap nozzle. Attempts have therefore already been made to blow through the central tube of the burner, into which a central ignition burner can be inserted, particulate secondary fuel which then emerges, however, as a compact jet, in which the particular fuel fans out a little, flies too far in the rotary tubular furnace, forms too long a flame and does not bum out or bums out too late. Furthermore, attempts have been made to set the injected secondary fuels in rotation at the mouth of the burner lance, the result of this being that, in particular, the large secondary fuel particles of high specific gravity are thrown onto the periphery, in any event are thrown out of the flame cone, instead of burning in the flame."}
{"text": "Traditional brushless canned motor pumps have a pump section and a motor section. The motor section drives the pump section. The pump section includes an impeller having blades which rotate inside a casing. The impeller pumps fluid from a pump inlet to a pump outlet. The impeller is normally of the closed type and is coupled to one end of a motor shaft that extends from the motor section into the pump section where it affixes to an end of the impeller.\nThe motor section includes an electric motor having a stator and a rotor. The rotor is unitarily formed with the motor shaft inside the stator. With brushless DC motors, the rotor is actuated by electromagnetic fields that are generated by current flowing through windings of the stator. A plurality of magnets are coupled to the rotor. During pump operation, the rotor shaft transmits torque, which is created by the generation of the electromagnetic fields with regard to the rotor's magnets, from the motor section to the pump section where the fluid is pumped.\nBecause the rotor and stator are immersed, they must be isolated to prevent corrosive attack and electrical failure. The rotor is submerged in the fluid being pumped and is therefore “canned” or sealed to isolate the motor parts from contact with the fluid. The stator is also “canned” or sealed to isolate it from the fluid being pumped. Mechanical contact bearings may be submerged in system fluid and are, therefore, continually lubricated. The bearings support the impeller and/or the motor shaft. A portion of the pumped fluid can be allowed to recirculate through the motor section to cool the motor parts and lubricate the bearings.\nSeals and bearings are prone to failure due to continuous mechanical wear during operation of the pump. Mechanical rub between the stator and the rotor can generate particles. Interacting forces between the rotor and the stator in fluid seals and hydrodynamic behavior of journal bearings can lead to self-excited vibrations which may ultimately damage or even destroy rotating machinery. The bearings are also prone to failure. Lubricants can be rendered ineffective due to particulate contamination of the lubricant, which could adversely affect pump operation. Lubricants can also dissolve in the fluid being pumped and contaminate the fluid. Bearings operating in a contaminated lubricant exhibit a higher initial rate of wear than those not running in a contaminated lubricant. The bearings and the seals may be particularly susceptible to failure when in contact with certain chemistry. Alternatively, the bearings may damage the fluid being pumped.\nWhat is needed is an improved brushless compact canned pump assembly design that substantially reduces particle generation and contamination, while rotating at high speeds and operating at supercritical temperatures and pressures."}
{"text": "Congestive heart disease is a progressive and debilitating illness characterized by a progressive enlargement of the heart. As the heart enlarges, the heart must perform an increasing amount of work to pump blood each heartbeat needed for metabolism. In time, the heart becomes so enlarged the heart cannot adequately supply blood. An afflicted patient is fatigued, unable to perform even simple exerting tasks and experiences pain and discomfort. Further, as the heart enlarges, the internal heart valves may not adequately close. This may impair the function of the valves and further reduce the heart's ability to supply blood. Importantly, there is no cure for congestive heart disease.\nCongestive heart failure has an enormous societal impact. In the United States alone, about five million people suffer from the disease. Alarmingly, congestive heart failure is one of the most rapidly accelerating diseases (about 400,000 new patients in the United States each year). Economic costs of the disease have been estimated at $38 billion annually. Not surprising, substantial effort has been made to find treatments for congestive heart disease.\nCardiomyoplasty is one potential treatment for moderate stage congestive heart disease. In this procedure, the latissimus dorsi muscle (taken from the patient's shoulder) is wrapped around the heart and chronically paced to assist contraction of the heart during systole.\nOne study speculates that an artificial elastic “sock” could be used in place of the latissimus dorsi in adynamic cardiomyoplasty. Kass et al., “Reverse Remodeling from Cardiomyoplasty in Human Heart Failure,” Circulation 91:9, May 1, 1995. Another study demonstrates that Bard Marlex sheets can be used to wrap the heart as a substitute to the latissimus dorsi in adynamic cardiomyoplasty. Oh et al., “The Effects of Prosthetic Cardiac Binding and Adynamic Cardiomyoplasty in a Model of Dilated Cardiomyoplasty,” Journal of Thoracic Cardiovascular Surgery, 116:1, July 1998. German Utility Model Patent Application DE 295 17 393 U1 describes a pericardium prosthesis made from a biocompatible, non-expansible material, or at least hardly expansible material that surrounds the heart to prevent over-expansion of the heart wall. PCT application WO 98/58598 describes an elastic pouch for at least partially enveloping a heart. Commonly assigned U.S. Pat. No. 5,702,343 to Alferness dated Dec. 30, 1997 teaches a jacket to constrain cardiac expansion during diastole. Other teachings include those of commonly assigned U.S. Pat. No. 6,123,662 and those of U.S. Pat. Application Publication No. US 2002/0019580."}
{"text": "Radio frequency (RF) signals often must be switched between two destinations, such as when switching an RF signal between a first antenna and a second antenna. Switches that support this configuration are classified as single pole, double throw (SPDT) switches.\nSPDT switches known in the art are either solid-state devices or mechanical relays. Mechanical switches are generally quite large, compared to other RF components, and consume significant amounts of power. In communications applications, switches are often designed with semiconductor elements such as transistors or pin diodes. However, for high frequency RF signals in the gigahertz range and above, these devices could suffer from several shortcomings. For example, the device has high insertion loss, which is the loss across the switch when the switch is closed. On the other hand, the device has low isolation, which means signal ‘leakage’ across the switch when the switch is open. In addition, the insertion losses and isolation values for these switches varies depending on the frequency of the signal passing through the switches, which cause unevenness in frequency response and is unsuitable for broadband wireless communications. Isolation between the two output ports of the SPDT switch is of another concern, since coupling of the signal from one output port to the other output port limits the effectiveness of the switch as a dual output port device. Without proper circuitry to overcome the above-mentioned characteristics, semiconductor transistors and pin diodes would not function properly as switches in microwave applications.\nFIG. 1A shows a basic prior art SPDT circuit, while FIG. 2A shows an improved prior art SPDT circuit with a parallel inductor. In FIG. 1A, a first switching device T101 such as a FET or PHEMT (Pseudomorphic High Electron Mobility Transistor) is positioned between the first port P1 and the second port P2. Correspondingly, a second switching device T102 is positioned between the second port P2 and the third port P3.\nA simplified equivalent circuit with transistor off can be represented by a parasitic capacitor between the drain terminal and the source terminal. The insertion loss of the SPDT switch is determined by the resistance value between drain and source of each transistor. As shown in FIG. 1B, when a semiconductor switching device is turned on, the device can be modeled with a resistor R. When a switching device is turned off, the device can be modeled with a capacitor C. Thus, when the first switching device is turned on and the second switching device is turned off, an equivalent circuit between P2 port and P3 port can be modeled as a resistance in series with a capacitance, with P1 port at the junction between the resistance and capacitance. The resistance results in loss of signal strength, which is commonly known as insertion loss. In addition, the signal can leak through the capacitance of the “off state” switching device, especially at high frequency. For instance, as shown in FIGS. 1C and 1D, using R=4 ohms and C=0.4 pF, there is a −1.7 dB insertion loss in transmission at 5 GHz between P1 and P2 and −6.5 dB of isolation at 5 GHz between P1 and P3.\nTo reduce the insertion loss, inductors can be added in parallel to the switching devices. The inductance and the capacitance of the switching device becomes a LC resonator with high impedance at the resonating frequency. For example, two inductors L are added in parallel to the source and the drain ends or terminals of the switching device T201 and T202 of FIG. 2A. When the first switching device T201 is turned on and the second switching device T202 is turned off, the circuit of FIG. 2A can be modeled by an equivalent circuit as a resistor R in a parallel with the inductor L and a capacitor C in parallel with inductor L, as shown in FIG. 2B. Assuming that L=2.5 nH, the insertion loss in transmission between P1 and P2 has been improved to −0.5 dB at 5 GHz, as shown in FIG. 2C. Moreover, the isolation between P1 and P3 has been improved to −27 dB at 5 GHz as shown in FIG. 2D.\nU.S. Pat. No. 5,774,792 uses such an approach where an SPDT switch includes a plurality of FETs, the FET on the receiver side through which a received signal passes and the shunt FET on the transmitter side are each formed of series-connected FETs, and a capacitor is connected between the first gate and the source and between the second gate and the drain. An inductance is connected in parallel with a series connection of FETs. Similarly, U.S. Pat. No. 6,693,498 uses a plurality of FETs in the SPDT switch and a first inductor connected in parallel to the first FET, and a second inductor connected in parallel to the second FET. The SPDT switch may further include a first inductor connected in series to the first FET, a first capacitor connected in parallel to the series connection of the first FET and the first inductor, a second inductor connected in series to the second FET, and a second capacitor connected in parallel to the series connection of the second FET and the second inductor.\nIn integrated circuits, size and space constraints are of primary importance. One problem with adding multiple inductors or other passive RF components to microcircuits is that the increase of the die size and thus increasing the cost of a switch."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for authentication of a user on the basis of his/her voice profile, particularly with the use of a mobile phone. Furthermore, the invention relates to a subscriber's identification module or SIM card for controlling such a method.\n2. The Prior Art\nVoice portals are gaining increasing significance. Because of the further development of speech recognition technology and because of the competition on this sector, call center services are being increasingly automated. Banks offer account status information, stock trading, and transfers via voice portal. Telecommunications providers make the Internet accessible by voice, via so-called Voice XML (extensible markup language).\nFor many services, the user must first prove his/her identity, i.e. must authenticate himself/herself. Conventional authentication by way of a secret number or a password, usually referred to as a PIN (Personal Identification Number), creates problems. These problems are created particularly if the service is used infrequently, if the user forgets his/her PIN, and consequently does not use the service, or if the user does not like to use the service, because of complicated PIN re-set procedures. Furthermore, there is the not insignificant security risk that an unauthorized third party will spy out the PIN and will gain access to the services of the aforementioned type, in place of the authorized user. Therefore voice service providers are beginning to authenticate the user by way of his/her voice profile. This form of authentication, referred to as voice authentication hereinafter, has obvious advantages with regard to the aspect of user acceptance. The user no longer needs to remember any secret data. He/she is guided through a short authentication dialog in the conversation with the computer that he/she has to conduct in any case. Subsequently the user receives access to the appropriate service, depending on the result of the voice check.\nFundamentally, there is the possibility of using the mobile phone terminal (handset), something that a user generally has, for utilizing the voice services and for a related voice authentication. However, the mobile phone does not offer a uniform voice transmission quality. In the case of weak field intensity, voices can be distorted and interference noises can be generated. Voice recognition, i.e. a check of characteristic voice characteristics for authentication, reacts to this distortion and noise more sensitively than speech recognition, i.e. the recognition of spoken information. Field tests show a relatively high so-called False Rejection Rate (false rejection of an authentic user) in mobile phone networks, and this high False Rejection Rate makes the use of mobile phone terminals for voice authentication questionable, in general.\nAside from the usual approach of obtaining a voice profile of the user on the provider side, from the transmitted voice signal, and comparing it with the reference voice profile of the user, another approach is also known. In this approach, the reference voice profile is stored on the subscriber's identification module or SIM card of the mobile phone terminal and the check algorithm is allowed to run on this SIM card. However, in view of the common memory space on SIM cards, great losses with regard to the quality of the reference voice profile and the check algorithm must be accepted. These losses in turn lead to high error rates and therefore make practical use fundamentally questionable."}
{"text": "(1) Field of the Invention\nThe present invention relates to a method for making a plasma-enhanced chemical vapor deposited (PECVD) silicon oxide/silicon nitride multilayer (stacked layer) as a passivation layer on semiconductor integrated circuits, and more particularly relates to a method for making an insulating/passivating layer over interconnecting metal lines on multilevel interconnections for integrated circuits.\n(2) Description of the Prior Art\nMultilevels of metal interconnections are used on Ultra Large Scale Integration (ULSI) integrated circuits to wire-up the discrete semiconductor devices on the semi-conductor chips. The different levels of interconnections are separated by layers of insulating material, such as PECVD silicon oxide (SiO.sub.2) and silicon nitride (Si.sub.3 N.sub.4) layers. The silicon nitride also serves as a barrier layer or passivation layer to prevent contamination from moisture and other corrosive chemicals, and to prevent ionic contamination from light alkaline metal ions, such as sodium (Na) and potassium (K) which can degrade the electrical characteristics of the underlying semiconductor devices built in and on the substrate. Via holes etched in the interposed insulating/passivation layers are used to connect one level of metal to the next. Typically the insulation and passivation layers require low-temperature deposition processing (&lt;400.degree. C.) because the metal lines which are usually formed from Al/Cu alloys have a low melting temperature (about 660.degree. C.).\nUnfortunately, as the metal lines that form the interconnections are made closer together to satisfy the higher packing density for ULSI integrated circuits, the aspect ratio (the width of the spacings between metal lines to the metal thickness) increases. This can result in the formation of voids (keyholes) in the nonconformal insulating or passivation layer between the metal lines when the passivation is deposited next. Also because the silicon nitride has a higher stress than the silicon oxide, cracking of the silicon nitride layer is prevalent at the bottom corners in the recesses between the metal lines. The voids are particularly prevalent when the passivation layer, usually a silicon nitride, is deposited at low temperatures, such as by PECVD or high-density plasma CVD (HDP-CVD) that results in the nonconformal layer.\nThe void and cracking problems are best understood by referring to the prior art in FIG. 1, which is a schematic cross-sectional view of a portion of a partially completed semiconductor substrate 18, such as a silicon substrate with an insulating layer 19 on the surface. The underlying semiconductor devices, such as field effect transistors (FETs) and electrical interconnections are not depicted to simplify the drawing and discussion. Two closely spaced interconnecting metal lines 20 are shown on which is deposited first a thin PECVD silicon oxide layer 22 that acts as a stress-release layer for the Si.sub.3 N.sub.4 layer 24 which is deposited next. As is depicted in FIG. 1, the deposition flux of the deposited Si.sub.3 N.sub.4 layer is higher at the top corners of the closely spaced metal lines than in the recess between them. This results in the formation of voids 30 having a keyhole shape. If the Si.sub.3 N.sub.4 is deposited to a thickness 24', as depicted by the broken line 24', then the voids can be partially closed. During the next photo-resist masking step to form the via holes in the passivation layer to the metal lines elsewhere on the substrate 18 (not shown), residual photoresist can be trapped in the voids 30 in the Si.sub.3 N.sub.4 layer 24. The incomplete removal of the residual photoresist then results in contamination and corrosion of the subsequent metal layer when the trapped photoresist outgases during the subsequent metal or alloy deposition step. Concurrently, stress-induced cracks 32 are formed in the silicon nitride layer 24 at the bottom edges of the metal lines 20, as is also depicted in FIG. 1. This is a particularly serious problem for 0.38 micrometer DRAM devices (0.38 um FET channel lengths) when the minimum spacings between the metal lines at the second metal level (M2) of interconnections are less than 0.575 um. The passivation layer typically fails the pinhole test when the total thickness of the silicon oxide layer 22 and the silicon nitride layer 24 is less than 8000 Angstroms, but thicker passivation layers can present a keyhole problem.\nOne method for depositing a passivation layer is described by Lindenfelser in U.S. Pat. No. 4,692,786, in which a three-layer sandwich passivation coating is deposited directly on the silicon substrate to prevent ion contamination during the fabrication of bipolar transistors, but he does not address the above keyhole or Si.sub.3 N.sub.4 cracking problems. Ang et al., U.S. Pat. No. 4,707,721, describe a method for making a passivated dual-gate system using a titanium nitride which is a good electrical conductor, but do not address the keyhole or cracking problem in a silicon nitride passivation layer.\nThere is still a strong need in the semiconductor industry to provide a silicon nitride passivation layer over closely spaced metal lines that eliminates the above keyholes and cracking problems on future ULSI and VLSI circuits."}
{"text": "1. Field of the Invention\nThe present invention relates to an exposure technique for exposing a substrate via a mask and a projection optical system, which is suitable for exposing a long sheet-shaped photosensitive object wound in a roll shape. In addition, the present invention relates to a fabricating technique for fabricating (manufacturing) a device using such an exposure technique.\n2. Description of Related Art\nFor example, in an exposure apparatus that is used for fabricating an element such as a semiconductor element or a liquid crystal display element, a typical exposure target (exposure object) has conventionally been a flat plate-shaped object which has high rigidity, such as a glass substrate or a semiconductor wafer which have been coated with photoresist. In recent years, in order to efficiently fabricate a device having a large area, as an exposure target, a long sheet-shaped member which is flexible and can be stored by being wound up into a roll shape has been used. When exposing such a long sheet-shaped member, as a conventional example, a part of the sheet-shaped member is exposed via a mask and a projection optical system is subjected to one-shot exposure (still exposure), the sheet-shaped member is then suctioned and moved, for example, by an exposure table, and then the suction is released and only the exposure table is returned to the initial position. The above described operation is repeatedly carried out to intermittently move the sheet-shaped member from a supply roller to a winding-up roller (for example, see Japanese Patent Application Publication No. 2007-114385).\nAccording to the conventional exposure method of the sheet-shaped member, one-shot exposure via a mask and movement of the sheet-shaped member are alternately repeated. Therefore, exposure efficiency is low, and, for example, it takes a long time to expose a sheet-shaped member which is wound up into one roll.\nIn contrast, a method is also considered in which the sheet-shaped member is scan-exposed while a mask and the sheet-shaped member are moved synchronously. However, in the case of mere scan exposure, for example, if the sheet-shaped member is exposed in an outward path of the mask, then it is difficult to expose the sheet-shaped member in a return path of the mask. For this reason, substantially half of the reciprocation time of the mask becomes a period at which exposure is impossible, and exposure efficiency cannot be high."}
{"text": "The increased speed and capability of computers and other electronic devices requires better performance from the integrated circuits that make up a device. One way to make the integrated circuits faster is to reduce the size of the transistors that make up the device. However, as transistors are made smaller and faster, delays through the connections between the transistors becomes greater in relation to the speed of the transistor.\nAn alternative technique to speed up integrated circuits is to use alternative semiconductors. For example, silicon-on-insulator (SOI) technology provides a 25-35% performance increase over equivalent CMOS technologies. SOI refers to placing a thin layer of silicon on top of an insulator such as silicon oxide or glass. The transistors would then be built on this thin layer of SOI. The SOI layer reduces the capacitance of the transistors so that they operate faster.\nFIG. 1 illustrates a typical prior art SOI semiconductor. The transistor is formed in the silicon layer 101 that is over the insulator 102. The insulator is formed on top of the substrate 103. Within the silicon layer 101, the drain/source regions 105 and 106 are formed. The gate 107 is formed above the partially depleted channel 109. A floating body 110 is within the depleted region 112 and results from the partial depletion.\nSOI technology, however, imposes significant technical challenges. The silicon film used for SOI transistors must be perfect crystalline silicon. The insulator layer, however, is not crystalline. It is very difficult to make perfect crystalline silicon-on-oxide or silicon with other insulators since the insulator layer's crystalline properties are so different from the pure silicon. If perfect crystalline silicon is not obtained, defects will find their way onto the SOI film. This degrades the transistor performance.\nAdditionally, floating body effects in partially depleted CMOS devices using SOI technology are undesirable in many logic and memory applications. The floating bodies cause threshold voltages and switching speeds to be variable and complex functions of the switching history of a particular logic gate. In dynamic logic and DRAM memories, the floating bodies cause excess charge leakage and short retention times that can cause data loss. In conventional flash memories and NROM devices, the floating bodies cause reduced erase fields and slower erase times.\nFor the reasons stated above, and for other reasons stated below which will become apparent to those skilled in the art upon reading and understanding the present specification, there is a need in the art for a way to eliminate floating body effects in CMOS devices incorporating SOI technology."}
{"text": "In a system, such as a basic business system for doing basic business and the like, since the scale of generated/stored data is large, the efficient execution of a process, such as data totaling, profiling, cleansing, add-up and the like, is required.\nIn relation to the above-described technology, a technology for automatically obtaining necessary data from a file (data) stored by a plurality of applications, using the statistic hydra of a TRIE structure is known.\nThe “statistic hydra” described here is the prior art different from “statistic hydra” described later as the best Mode for carrying out the invention.\nFor example, when a test result data 1200 illustrated in FIG. 1 is classified into “very good”, “good”, “fairly good” and “bad” according to an evaluation table 1300 illustrated in FIG. 2 and in order to obtain their distribution and its average, the following processes are needed.\n(1) A statistic hydra 1400 illustrated in FIG. 3 is generated from the test result data 1200 and a primary totaling is applied to each score. As a result, a totaling result 1401 for each “score” is obtained.\n(2) The data of the totaling result 1401 is classified into respective “score ranges” of the evaluation table 1300, for example, using the following retrieval expressions. As a result, as illustrated in FIG. 4, the totaling results 1501, 1502, 1503 and 1504 for each “score range” of the evaluation table 1300 are obtained.\nTotal data group retrieval expression of “very good”: Score≧45 AND score≦50\nTotal data group retrieval expression of “good”:\nScore≧40 AND score<45\nTotal data group retrieval expression of “fairly good”: Score≧30 AND score≦40\nTotal data group retrieval expression of “bad”: Score≧0 AND score<30\n(3) Since a plurality of totaling results is obtained in the totaling results 1503 and 1504, the secondary totaling is executed in order to obtain a totaling result for each “score range”. As a result, a desired totaling result 1600 illustrated in FIG. 5 is obtained.    Patent document 1: Japanese Laid-open Patent Publication No. 2006-171800"}
{"text": "1. Field of the invention\nThe invention relates to a semiconductor device and a method of manufacturing the same, specifically relates to a semiconductor device with a semiconductor chip encased in a hollow ceramic package and a method of manufacturing the same.\n2. Description of the related art\nA conventional method for encasing a semiconductor chip in a hollow ceramic package is disclosed in the reference [JP 2006-179740A]. According to JP 2006-179740A, a ceramic sheet, which is divided into a plurality of individual headers (containers) later, and a metal plate, which is divided into a plurality of individual lids for the headers later, are used for the ceramic package. The ceramic sheet includes a plurality of headers disposed in a matrix, and each header includes a cavity, which has an opening located at an upper surface of the ceramic sheet, for placing a semiconductor chip therein. The ceramic sheet also includes grooves formed at a back surface, which is opposite to the upper surface, for defining the headers. In order to manufacture a semiconductor device with a semiconductor chip encased in a hollow ceramic package, the semiconductor chip is placed in each cavity, and then is connected to the header by bonding wires. Then, the metal plate is placed on the ceramic sheet by sandwiching a sealing sheet. After that, the metal plate is fixed to the ceramic plate by using a roll electrode, which melts down the sealing sheet by electric power applied to the roll electrode. Then, after only the metal plate is diced, the ceramic sheet is divided into individual semiconductor devices in a so-called chocolate-break manner by an external force along the grooves formed therein.\nHowever, according to the conventional method, since the surface of the metal plate on which the roller electrode is contacted is a flat surface, the contact area of the roller electrode and the metal plate is limited so that a high voltage is required to melt-down the sealing sheet for fixing the metal plate to the ceramic sheet. Further, it is not easy to divide the metal plate only by the dicing method. Moreover, two dividing steps are required in the conventional method. Namely, the first step is a dicing process for dicing the metal plate only, and the second step is a chocolate-breaking process for breaking a ceramic sheet in order to divide into the individual semiconductor devices. The complex process makes the manufacturing process complex, decreases the productivity, and increase the cost for production."}
{"text": "1. Field of the Invention\nThe present invention relates to a delay circuit, a testing apparatus, and a capacitor. More particularly, the present invention relates to a delay circuit which generates a desired delay time by changing the junction capacitance of a field effect transistor.\n2. Description of the Related Art\nFIG. 1 shows a conventional delay circuit 300. The conventional delay circuit 300 has a first buffer 302 which shapes the wave form of an input signal and then outputs the resultant shaped signal, a path 306 through which the output signal transmits, a first capacitor 312 which adds capacitance C to the path 306, a second capacitor 314 which adds capacitance C′ to the path 306, a first switching device 308 which electrically connects or disconnects the path 306 with the first capacitor 312, a second switching device 310 which electrically connects or disconnects the path 306 with the second capacitor 314, and a second buffer 304 which shapes the wave form of the signal that has transmitted through the path 306 and outputs the resultant shaped signal. A control unit not shown in the drawing controls the switching devices 308 and 310 so as to change the capacitance added to the path 306. In this way, the control unit not shown in the drawing delays the signal that transmits the path 306 by a desired length of time.\nThe conventional delay circuit 300 achieves a fine delay resolution by selectively adding either the capacitance C or the capacitance C′ which differs slightly from the capacitance C. However, in the conventional delay circuit 300, the channel capacitance of the first switching device 308 differs from that of the second switching device 310, and the wire capacitance of the wire which connects the first capacitor 312 with the path 306 differs from the wire capacitance of the wire which connects the second capacitor 314 with the path 306. These capacitance differences influence the capacitance added to the path 306. As a result, the desired fine delay resolution which is designed to be achieved by utilizing the fine difference between the capacitance C and the capacitance C′ has been very difficult."}
{"text": "1. Field of the Invention\nThis invention relates to a semiconductor laser device and a manufacturing method thereof. In particular, it relates to a semiconductor laser device used to process optical information and a manufacturing method thereof.\n2. Description of the Related Art\nA semiconductor laser device for optical information processing has conventionally employed a gain guided structure using a GaAs current blocking layer. Recently, there has been, however, developed a semiconductor laser device which employs a real refractive index guided structure using an AlInP layer as a current blocking layer to reduce an operating current.\nA real refractive index guided structure may reduce an optical absorption loss in a current blocking layer resulting in not only a reduced threshold current but also an improved luminous efficiency, therefore a reduced operating current.\nThis technical trend has been driven for developing semiconductor laser device shaving a higher output. Conventional optical information processing involves only reading as in, for example, DVD-ROM, which does not require very high output. Recent optical information processing involves, however, not simply reading but also writing on a recording medium as in, for example, DVD-RW or DVD-R, which necessarily requires a higher output. It has been, therefore, required that an internal loss is minimized to reduce an operating current for improving temperature properties of the semiconductor laser device and thus reliability under a high output.\nFIG. 14 is a cross-sectional view of a conventional SAS (Self-Aligned Structure) type of red semiconductor laser diode (hereinafter, referred to as a xe2x80x9cred LDxe2x80x9d) described in Electronics Letters, Vol. 33, No. 14 (1997), pp.1223-5.\nIn FIG. 14, reference numeral 100 denotes a red LD, 102 an n-type GaAs substrate (hereinafter, n-type and p-type are denoted as xe2x80x9cn-xe2x80x9d and xe2x80x9cp-xe2x80x9d, respectively), 104 an n-GaAs buffer layer, 106 a lower clad layer made of n-(Al0.7Ga0.3)0.5In0.5P, and 108 an active layer of an MQW structure made of GaInP/AlGaInP where GaInP is a material for a well layer and AlGaInP is a material for a barrier layer.\nIn this figure, reference numeral 110 denotes a first upper clad layer made of p-(Al0.7Ga0.3)0.5In0.5P, 112 a current blocking layer made of n-AlInP, 114 a stripe-shaped opening to be a current channel in the current blocking layer 112, 116 a second upper clad layer made of p-(Al0.7Ga0.3)0.5In0.5P, 118 a p-GaAs contact layer, 120 a p-electrode, and 122 an n-electrode.\nThere will be described a process for manufacturing this semiconductor laser device 100.\nFIGS. 15, 16 and 17 are cross-sectional views of a conventional red LD in individual manufacturing steps.\nFirst, on an n-GaAs substrate 102 are sequentially deposited an n-GaAs layer to be a buffer layer 104, an n-(Al0.7Ga0.3)0.5In0.5 layer to be a lower clad layer 106, a GaInP/AlGaInP MQW layer to be an active layer 108, a p-(Al0.7Ga0.3)0.5In0.5P layer to be a first upper clad layer 110 and an n-AlInP layer to be a current blocking layer 112, by primary epitaxial growth based on crystal growth such as MOCVD. For dopants, silicon is used as an n-type dopant while zinc is used as a p-type dopant. The result of this step is shown in FIG. 15.\nThen, a resist pattern 126 is formed on the surface of the n-AlInp layer to be a current blocking layer 112 by a photolithographic process, and a stripe-shaped opening 114 to be a current path is formed in the n-AlInP layer to be the current blocking layer 112 by wet etching. The result of this step is shown in FIG. 16.\nAfter removing the resist pattern 126, a p-(Al0.7Ga0.3)0.5In0.5P layer to be a second upper clad layer 116 is formed on the p-(Al0.7Ga0.3)0.5In0.5P layer to be a first upper clad layer 110 facing the opening 114 and the n-AlInP layer to be a current blocking layer 112 by secondary epitaxial growth based on crystal growth such as MOCVD. The result of this step is shown in FIG. 17.\nThen, a p-GaAs layer to be a contact layer 118 is formed on the p-(Al0.7Ga0.3)0.5In0.5layer to be a second upper clad layer 116.\nIn this process, crystal growth temperature is about 650xc2x0 C. to 750xc2x0 C. A crystal growth temperature as low as possible is used to prevent the p-type dopant, Zn, from diffusing from the p-(Al0.7Ga0.3)0.5In0.5P layer, as the first upper clad layer 110, into the MQW layer, the active layer 108, to the maximum extent possible.\nThen, a p-electrode 120 and an n-electrode 122 are formed on the surface of the p-GaAs layer to be a contact layer I 10 and on the rear surface of the n-GaAs substrate 102, respectively.\nA conventional red LD 100 has a configuration as described above. When forming the p-(Al0.7Ga0.3)0.5In0.5P layer for the second upper clad layer 116 on the n-AlInP layer for the current blocking layer 112 in the manufacturing process for the red LD 100 as illustrated in FIG. 17, lattice defects frequently develop on the surface facing the opening 114 in the n-AlInP layer for the current blocking layer 112, leading to an increase in internal loss of light, deterioration in temperature properties, and poor reliability of the red LD 100.\nA technique for preventing lattice defects in crystal growth has been described in Proceedings of the Tenth International Conference on Metal organic Vapor Phase Epitaxy (2000), p. 82. In the report, an (Al0.7Ga0.3)0.51In0.49P layer is deposited on a GaAs substrate which is a (100) facet misoriented by 10xc2x0 toward [011] direction. Then, on the layer is formed an Al0.51In0.49P layer having a grooved structure whose side wall is a (111) A facet, and with Ga0.51In0.49P as a marker sandwiched in between, an (Al0.7Ga0.3)0.51In0.49P layer is formed above the (Al0.7Ga0.3)0.51In0.49P layer parallel to the GaAs substrate exposed in the bottom of the grooved structure and above the Al0.51In0.49P layer having a (111) A facet, during which development of lattice defects is studied using the then substrate temperature as a parameter.\nAccording to the report, crystal growth was caused at substrate temperatures of 720xc2x0 C., 760xc2x0 C. and 800xc2x0 C. It was found that lattice defects developed in a crystal layer growing on a (111) A facet at a substrate temperature of 720xc2x0 C. or 760xc2x0 C., while crystal growth at a substrate temperature of 800xc2x0 C. reduced lattice defects in the (Al0.7Ga0.3)0.51In0.49P layer on a (111) A facet.\nHowever, for a red LD, a crystal growth temperature of 800xc2x0 C. may cause diffusion of the p-type dopant Zn from the p-(Al0.7Ga0.3)0.5In0.5P layer as a first upper clad layer 110 to the MQW layer to be an active layer 108, leading to deterioration in temperature properties or reliability in current-optical output performance.\nBesides the prior art described above, JP-B 2842465 has disclosed an SAS type semiconductor laser where on the surface of a current blocking layer made of AlGaAs material having a stripe-shaped opening is deposited a protective layer made of an AlGaAs material with small aluminum content, on which a p-AlGaAs material is deposited as a p-clad layer, but has not described that on a current blocking layer made of an AlInP material are formed a capping layer made of a GaInP material and a p-(Al0.7Ga0.3)0.5In0.5P layer.\nThe present invention has been made to solve the above problem in the art, and an objective of this invention is to provide a reliable semiconductor laser device exhibiting a reduced threshold current with less deterioration in temperature properties in current-optical output performance.\nA semiconductor laser device according to the present invention comprises: a semiconductor substrate of a first conductivity type; a first clad layer of a first conductivity type made of a III-V group compound semiconductor disposed on the semiconductor substrate; an active layer made of a III-V group compound semiconductor having a smaller band gap than the first clad layer, disposed on the first clad layer; a first second-clad layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer, disposed on the active layer; a current blocking layer of a first conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer, disposed on the first second-clad layer and having a stripe-shaped opening to be a current path; a buffer layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer, disposed on the surface of the current blocking layer facing the opening; and a second second-clad layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer, disposed on the first second-clad layer facing the opening and the current blocking layer via the buffer layer.\nAccordingly, a semiconductor laser device according to the present invention is advantageous in that lattice defects in the second second-clad layer disposed on the surface of the current blocking layer facing the opening via the buffer layer may reduce and it can be prevented from a second conductivity type of dopant diffusing from the first second-clad layer to the active layer. Therefore the construction according to the present invention makes it possible to reduce deterioration in temperature properties in current-optical output performance, and consequently to improve reliability of the semiconductor laser device.\nAnother objective of this invention is to lead to improvement in reliability of a semiconductor laser device with a red LD.\nA semiconductor laser device according to the present invention comprises: a semiconductor substrate of a first conductivity type; a first clad layer of a first conductivity type made of a III-V group compound semiconductor disposed on the semiconductor substrate; an active layer made of a III-V group compound semiconductor having a smaller band gap than the first clad layer, disposed on the first clad layer; a first second-clad layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer, disposed on the active layer; a current blocking layer of a first conductivity type made of an AlInP material having a larger band gap than the active layer, disposed on the first second-clad layer and having a stripe-shaped opening to be a current path; a protective layer made of a GaInP material disposed on the surface of the current blocking layer except the surface facing the opening; and a second second-clad layer of a second conductivity type made of an AlGaInP material having a larger band gap than the active layer, disposed on the current blocking layer via the protective layer and the first second-clad layer facing the opening.\nAccordingly, a semiconductor laser device according to the present invention is advantageous that because the current blocking layer made of the AlInP material is protected by the GaInP material not containing Al the second second-clad layer made of the AlGaInP material whose composition tends to be deviated is disposed on the crystal facet with a reduced amount of oxide film. Consequently reducing a risk of lattice defect formation leads to improvement in reliability of a semiconductor laser device with a red LD.\nA further objective of this invention is to provide a process for manufacturing a reliable semiconductor laser device exhibiting a reduced threshold current with less deterioration in temperature properties in current-optical output performance by simple steps.\nA process for manufacturing a semiconductor laser device according to the present invention includes the steps of: depositing a first clad layer of a first conductivity type made of a III-V group compound semiconductor, an active layer made of a III-V group compound semiconductor having a smaller band gap than the first clad layer, a first second-clad layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer, and a current blocking layer of a first conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer on the first second-clad layer in order, on a semiconductor substrate of a first conductivity type; forming a stripe-shaped opening penetrating the current blocking layer; forming a buffer layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer on the surface of the current blocking layer facing the opening; and forming a second second-clad layer of a second conductivity type made of a III-V group compound semiconductor having a larger band gap than the active layer on the first second-clad layer facing the opening and the current blocking layer via the buffer layer.\nAccordingly, a process for manufacturing a semiconductor laser device according to the present invention is advantageous that because the crystal of the second second-clad layer may grow with reduced lattice defects on the surface of the current blocking layer facing the opening via the buffer layer even at a common substrate temperature and such a common substrate temperature may reduce diffusion of a second conductivity type of dopant from the first second-clad layer to the active layer during crystal growth, a reliable semiconductor laser device with reduced deterioration in temperature properties in current-optical output performance can be manufactured by simple steps. Consequently, a semiconductor laser device exhibiting excellent laser properties can be produced with a low price.\nOther objects and advantages of the invention will become apparent from the detailed description given hereinafter. It should be understood, however, that the detailed description and specific embodiments are given by way of illustration only since various changes and modifications within the scope of the invention will become apparent to those skilled in the art from this detailed description."}
{"text": "DNA-topoisomerases are enzymes which are present in the nuclei of cells where they catalyze the breaking and rejoining of DNA strands, which control the topological state of DNA. Recent studies also suggest that topoisomerases are also involved in regulating template supercoiling during RNA transcription. There are two major classes of mammalian topoisomerases. DNA-topoisomerase-I catalyzes changes in the topological state of duplex DNA by performing transient single-strand breakage-union cycles. In contrast, mammalian topoisomerase II alters the topology of DNA by causing a transient enzyme bridged double-strand break, followed by strand passing and resealing. Mammalian topoisomerase II has been further classified as Type II α and Type II β. The antitumor activity associated with agents which are topoisomerase poisons is associated with their ability to stabilize the enzyme-DNA cleavable complex. This drug-induced stabilization of the enzyme-DNA cleavable complex effectively converts the enzyme into a cellular poison.\nSeveral antitumor agents in clinical use have potent activity as mammalian topoisomerase II poisons. These include adriamycin, actinomycin D, daunomycin, VP-16, and VM-26 (teniposide or epipodophyllotoxin). In contrast to the number of clinical and experimental drugs which act as topoisomerase II poisons, there are currently only a limited number of agents which have been identified as topoisomerase I poisons. Camptothecin and its structurally-related analogs are among the most extensively studied topoisomerase I poisons. Recently, bi- and terbenzimidazoles (Chen et al., Cancer Res. 1993, 53, 1332-1335; Sun et al., J. Med. Chem. 1995, 38, 3638-3644; Kim et al., J. Med. Chem. 1996, 39, 992-998), certain benzo[c]phenanthridine and protoberberine alkaloids and their synthetic analogs (Makhey et al., Med. Chem. Res. 1995, 5, 1-12; Janin et al., J. Med. Chem. 1975, 18, 708-713; Makhey et al., Bioorg. & Med. Chem. 1996, 4, 781-791), as well as the fungal metabolites, bulgarein (Fujii et al., J. Biol. Chem. 1993, 268, 13160-13165) and saintopin (Yamashita et al., Biochemistry 1991, 30, 5838-5845) and indolocarbazoles (Yamashita et al., Biochemistry 1992, 31, 12069-12075) have been identified as topoisomerase I poisons. Other topoisomerase poisons have been identified including certain benzo[i]phenanthridine and cinnoline compounds (see LaVoie et al., U.S. Pat. No. 6,140,328 (735.037WO1), and WO 01/32631(735.044WO1)). While these compounds are useful they are somewhat limited due to low solubility."}
{"text": "In magnetic and magneto-optical switching devices, in particular, used for recording purposes, information bits are stored as oppositely oriented magnetic domains representing the values “0” and “1” respectively. The conventional way of switching spins or magnetic domains is to apply an external magnetic field along the direction of spins, opposite to their initial orientation. This is a rather slow reversal process. Alternatively the external field can be applied perpendicular and not parallel to the spin orientation. This will trigger a much faster reversal process via a precessional motion of the spins around the external applied field. A plurality of magnetic and magneto-optical recording devices is being currently manufactured using such principles for reversing magnetization and thereby write and re-write the information.\nA new type of such magnetic recording devices designed to achieve ultra-high data storage densities, is the so called heat assisted magnetic recording (HAMR). Here a high anisotropy magnetic medium is heated by a laser beam such that the external applied field can still be efficient in reversing the magnetization. However, this approach does not simplify the magnetic recording process but instead involves more elements leading to high manufacturing costs of the magnetic devices as well as higher power consumption. Furthermore, by I. Tudosa et all., Nature 428, 831 (2004) and C. H. Back et al., Nature 428, 808 (2004) it has been demonstrated that the ultimate speed of magnetic switching triggered by a magnetic field is bounded in the picosecond time scale. Since the demand for ever increasing the density of data storage requires increasing of the magnetization switching speed too, new and faster ways of switching the magnetization other than by magnetic fields are desirable. Incorporated by reference, in Kimel et al. “Ultrafast non-thermal control of magnetization by instantaneous photomagnetic pulses”, Nature Letters, Published online 25 May 2006, further referred to as “Nature Letters article” non-thermal excitation was demonstrated for dielectric materials. Also incorporated by reference, in Hansteen et al. “Femtosecond photomagnetic switching of spins in Ferrimagnetic Garnet Films”, Physical Review Letters, 047402 (2005), further referred to as the PRL article, small angle magnetization precessions were demonstrated to have been successfully controlled. The materials in which optically induced effects were demonstrated are insulating materials."}
{"text": "1. Field\nThe present disclosure relates generally to non-destructive inspection and, in particular, to performing non-destructive inspection on a structure. Still more particularly, the present disclosure relates to a method and apparatus for detecting inconsistencies in a structure using a plurality of types of ultrasonic signals substantially simultaneously.\n2. Background\nIn manufacturing aircraft, vehicles, and other structures, inspection of parts used to form these structures is often performed to determine whether the parts will have desired parameters for a desired performance of the part. Additionally, the structures and parts are inspected as part of normal maintenance when the aircraft, vehicles, and other structures are in use.\nNon-destructive testing is commonly performed on these parts. Non-destructive testing is used to evaluate the properties of a part without altering the ability to use the part in service.\nUltrasound testing is a type of non-destructive testing. Ultrasound testing is often used to perform inspections on aircraft parts that include, or are comprised of, composite materials. Ultrasound testing involves transmitting acoustic waves through a test object, such as an aircraft part or structure.\nUltrasound testing is commonly performed using a transducer. The transducer is configured to send acoustic waves into a test object at a near side and detect a response to the acoustic waves. The response to these acoustic waves is analyzed to determine whether inconsistencies are present in the test object.\nDifferent types of inconsistencies may be identified using different types of acoustic waves. For example, inconsistencies on a far side of a structure may be identified using shear waves. A far side of the structure is a side of the structure opposite the near side into which the acoustic waves are transmitted. As another example, inconsistencies on a near side of the structure may be identified using surface waves. In conventional ultrasonic inspection, each type of acoustic wave may be detected on a separate pass of the structure. For example, a first pass of the structure may detect surface waves while a second pass of the structure may detect shear waves. However, performing multiple inspection passes of the structure may use an undesirable amount of time, energy, or resources. Therefore, it would be desirable to have a method and apparatus that take into account at least some of the issues discussed above, as well as other possible issues."}
{"text": "The present invention relates to a method and/or a device permitting, in connection with the control of electric motors, to extend to a certain degree, in the sense of improving, the stoppage safety function of such motors.\nThe invention is based on a monolithic integrated circuit which has long been known, corresponding for example to a U211B circuit chip of the type produced and sold by the Telefunken Company, which chip is used especially in modern phase control systems employed for controlling electric motors of any application.\nSo, it has been known for example (DE-OS 32, 22, 065) to equip an electric hand-held tool, especially a hand-held drilling or screwing machine, with a torque-limiting feature for the speed control, and to effect a speed adjustment in connection therewith. The torque-limiting feature acts in one sense of rotation only, for example in the screwing-in sense, and is automatically switched off in the other sense of rotation.\nIt has further been known to provide an electric motor of a hand-held tool intended for tightening nuts (U.S. Pat. No. 3,892,146) with an electric drive circuit which drives the motor, for the purpose of tightening of the nut, until the armature current of the motor exceeds a predetermined value; the drive circuit responds when this threshold value is exceeded and reverses the sense of rotation of the motor, thereby releasing the nut via an adapter.\nIn the case of another known power-driven screwing machine with torque-limiting feature (EU-OS 0,187,353; DE-PS 35,00,714) the electric motor is driven by means of a phase control comprising a setting potentiometer for the desired torque, with a usual triac being connected into the load circuit of the electric motor, whose angle of current flow is released by the phase control in response to the setting.\nAs has been mentioned before, such known phase control systems are always provided with incorporated speed safety functions which will respond when the electric motor is further excited, when it is already in a stopped or almost stopped condition under load.\nThe before-mentioned phase control chip model U211B from Telefunken, for example, is equipped with an input connection (pin 8) through which a frequency-to-voltage converter can be controlled. Due to a corresponding internal wiring system, each drive signal will result in recharging conditions at this input connection, which may be connected for example to a tachogenerator of the electric machine being monitored and driven, which recharging conditions are utilized for releasing the excitation of the triac.\nTo this end, an additional input connection (pin 18) is provided which offers a pulse-blocking feature which becomes effective when no recharging conditions are encountered any more in the frequency-to-voltage converter, as a result of incoming speed-responsive pulses. In this case, a capacitor connected to the input connection, corresponding to pin 18, is charged up, via a resistor, until a cut-off threshold is reached at which point the triac triggering pulses will be interrupted automatically by the control module corresponding to U211B.\nIn the case of such a circuit, which has been known as such and also in connection with driving systems for electric motors, considerable problems may be encountered in special applications, as a result of the tacho-monitoring and speed stoppage function, as follows:\nAs has been mentioned before, the output signal blocking function will not become active due to recharging processes in the f/v converter, the control circuit U211B being continuously supplied with speed-responsive pulses of any shape and magnitude. The ordinary and, insofar, welcome cut-off function (i.e. interruption of the triac triggering pulses) must however lead to faulty operation when such an electric motor is used for driving a power-driven screwing machine, which sometimes performs screwing operations at extremely high torques and with the aid of a greatly stepped-down intermediate transmission. Until the electric motor of such a power-driven screwing machine, in particular if the latter is designed for high and extremely high torques, is switched off by the speed-responsive stoppage monitoring function, the entire transmission train will be subjected to considerable stresses, including torsion of shafts, or the like, so that when the motor current is switched off via the triac of the phase control in the control module U211B, which is used here as control circuit by way of example, a certain return motion will occur in the transmission and motor area, and this will of course give rise again to the generation of new tachogenerator pulses by the existing tachogenerator or other means generating speed-responsive pulses.\nNow, it is exactly one of the properties of the before-mentioned control module, and practically of any other modern phase control circuit as well, that when speed-responsive pulses are received, certain processes, in the present case the recharging processes in the converter will be resumed and the values will drop below the cut-out threshold so that the electric motor will be driven again at full motor current. Consequently, the motor will be started again and run up as far as possible. But given the fact that the screw had been tightened before, the motor will be blocked and stopped again so that no speed-responsive pulses will be received any more and the circuit will switch the system off; but the high mechanical stresses, which are again encountered, result again in a return motion so that the whole system gets into pulsating operation. However, such a pulsating behavior is particularly undesirable for the power-driven screwing machines discussed here as it results in just that type of malfunction which was to be excluded by the torque-detection and setting features. In fact, all electronic control modules containing modern phase control systems are designed in such a way as to permit presetting of a torque so that when a given, maybe even very high torque has been selected power-driven screwing machines are capable of performing screwing operations perfectly, in particular smoothly, until the maximum torque has been reached, without there occurring any torque peaks, abrupt vibrations, or other disturbances. However, the malfunction just described, which is the result of the speed monitoring in the stopped condition, leads to quite the contrary condition so that it is no longer possible to preset defined torques, the mechanical kinematics in the system giving rise to sudden, abrupt torque peaks when the system starts running again, following a return motion, which torque peaks have particularly undesirable effects on the part to be screwed in or to be tightened, due to the high stepping-up ratio.\nNow, it is the object of the present invention to remedy these disadvantages and to extend the stoppage safety function, in particular for an electric motor used in power-driven screwing machines, in such a way that the electric motor will not start running up again, i.e. will not be fully switched on by the existing monolithic integrated phase control circuit, even when, after the system has been switched off for the first time, stresses produced in the system will give rise to new speed-responsive pulses, during the return motion, due to a reversing process."}
{"text": "This invention relates to a weaving heddle frame having lightweight frame staves, each frame having at least one heddle carrying rod located on one side thereof for taking up heddles with laterally open end loops.\nIn order to achieve a high working speed on modern weaving machines, it is essential that the heddle frame, which moves with the corresponding high frequency of strokes, be as lightweight as possible but nevertheless have the greatest possible bending strength. A high bending strength of the heddle frame, which is essential to keep the stress of the lined-up heddles within limits, could be achieved with frame staves of solid construction. However, in view of the required working speed and high frequency of strokes, it would be difficult to meet the lightweight requirement with such a solid construction. Heddle frames have therefore been made of lightweight construction including frame staves and heddle carrying rods for the heddles forming a single unit and being manufactured of the same material, the frame staves being of hollow construction. The heddle frames may be of light metal or of synthetic material which may be reinforced by fibers or webs.\nWith frame staves of the aforementioned type, forming a single unit with heddle carrying rods for the heddles, a sufficient resistance to wear and tear cannot however be achieved if, at each stroke of the heddle frame, the heddles with metallic end loops strike against the heddle carrying rod made of a softer material."}
{"text": "1 . Field of Invention\nThe present invention relates to windows and more particularly is directed to a thermal-break window fabricated of a plurality of extruded frame members designed, when interconnected, to form a substantially rectangular frame for encasing a glass panel.\n2. Description of the Prior Art\nWindows of various configurations have previously been made from a plurality of precut extruded sections that were fastened together by welding and/or by fasteners, e.g., screws, nuts and bolts or the like. All these windows have been suffered from the disadvantages of complex and costly fabrication and assembly techniques. Furthermore, all these prior windows possessed an undesirable characteristic, namely that of conducting heat or cold from the outside atmosphere through their metal frame members into the room, thus not being energy efficient."}
{"text": "As shown in FIG. 10, a speaker 100 that is a related-art common electromechanical acoustic transducer has a housing 101 having a sound port 101a at the center portion of the housing 101, and an essentially-plate-shaped frame 103 is attached to an interior surface 101b of the housing 101 through an rising portion 102 so as to generate an internal space 104.\nA cylindrical closed-end yoke 105 supported in correspondence with the internal space 104 is attached to the center portion of the frame 103. A magnet 106 is attached to the inside of a bottom 105a of the yoke 105, and a magnetic gap 107 is formed between an interior surface of the yoke 105 and the magnet 106.\nA plate 108 is attached to an apical surface of the magnet 106. A voice coil 109 is inserted in the magnetic gap 107. One end of the voice coil is attached to the bottom 105a of the yoke 105, and the other end of the voice coil 109 is attached to the center portion of a diaphragm 110.\nA protector 111 for protecting the diaphragm 110 is provided in front of the diaphragm 110, and first air holes 111a are provided in the protector 111, and second air holes 103a are provided in the frame 103.\nIn such a speaker 100, the bottom 105a of the yoke 105 is arranged so as to face from the interior surface 101b of the housing 101 toward a departing direction (a downward direction in FIG. 10). Since the yoke 105 is smaller in diameter than the frame 103, a space exists between the frame 103 and another electronic component M housed in the housing 101.\nIn order to achieve a superior acoustic characteristic, a necessity for assuring a sufficient space around the second air holes 103a for producing superior resonances of sound emitted from the second air holes 103a to the internal space 104 of the housing 101 has commonly been known (see; for instance, Patent Document 1).\nPatent Document 1: JP-A-2002-171596"}
{"text": "Modern techniques for constructing building walls include, for example, the two-by-four (2×4) framed construction, comprising wood or metal (such as, steel) members. These 2×4s are oriented vertically and connected at the top and bottom to similar members that are horizontally oriented. This structure is referred to in the relevant art as a “framed” wall. A sheet of building wall substrate or exterior sheathing wallboard, such as plywood sheathing, specially formulated gypsum wallboards, cement panels, or fiber cement panels are affixed to building frame.\nA water-resistive barrier is then typically applied to the exterior of the sheathing wallboard, with an external wall cladding or finish then being applied directly over the water-resistive barrier. Many materials may be used for the external wall finish such as brick, stucco, vinyl or aluminum siding, and/or wood. For example, stucco exterior finishes make up one of the most common ways of finishing exteriors of both residences and commercial buildings. Stucco finishes have limitations including, porosity, rigidity, freeze/thaw fractures, fungal and mildew formation, cracking, and compromised weatherability if not applied properly. However, because modern exterior finishes are watertight, any water that remains trapped behind these finishes does not readily evaporate. The trapped water behind the exterior finish can then soak into the substrata, for example, the sheathing wallboard and framing of the building. The water that soaks into the substrata and framing often causes severe damage to the building without any signs of damage appearing on the exterior of the building. These problems can exist regardless of the age of the building or the quality of construction.\nThe exterior sheathing wallboards such as gypsum (or other) wall-boards used as structures, are joined at edges by joint compounds. During construction, these exterior sheathing wallboards are exposed to outside environment, which, compared to the interior environment of a building is rather severe. The exterior sheathing wallboards are exposed to UV, rain, wind, heat, cold, and other physically abrasive impactors such as debris, insects and birds, particularly during construction work. Thus, these boards are specially formulated to have weatherability and particularly water resistance of 100 g/m2 or lower 2-hour water permeance or Cobb value. Just like interior drywall panels, exterior sheathing wallboard panels, when installed, have seams or joints between each panel that must be filled so that a smooth exterior finish may be applied. The joints between these exterior wall-boards or sheathing, therefore, become the weak link from where moisture can penetrate through the walls of the building and into the interior. Currently, the joint compound (or the water resistive barrier) used to seal the joints between the sheathing and to provide water resistance include a silyl-terminated-polyether (STPE) flashing membrane system. The Cobb value of such STPE systems is around 2 g/m2. Stated another way, the joint compound or the sealant that is currently used is about thirty times more water resistant than the gypsum wall-boards having the much larger exposed area. The STPE based systems are not water based and therefore have inherent issues relating to dealing with organic chemical solvents. Also, these are very expensive systems.\nFor the exterior or the interior wallboards, to achieve a smooth, visually appealing surface, the joints between boards, cracks, screw holes, and/or nail holes must be concealed. Conventional wallboard joint compounds for the interior are used to cover and finish gypsum wallboard joints, cornerbead, and screw or nail holes. Joint compounds can be spread over mesh or tape used to connect wallboards. It may also be used to patch and texture interior walls. The intrusion of water through wall spaces, either through prolonged direct contact or via high humidity, has a debilitating effect (mold and structural damage) on standard wall systems. The result of water seepage through joint compound to the studs on the other side of the wall ultimately has devastating structural and microbial implications for the wall system, first by absorption of the seeped water into the wood studs followed by their swelling and deformation (leading to expensive structural problems) and then, the creation of a fertile ground for rapid mold growth.\nSome performance requirements of exterior sheathing wallboards include (1) long-term weather-resistance; (2) mold resistance; and (3) water resistance. These performance characteristics are quite stringent, and it is to be expected that whichever material that is to be used to cover the joints between these sheathing boards (as well as the nail heads) should, at a minimum, display the same performance characteristics as that of the corresponding sheathing wallboards. It is problematic achieving all of these characteristics, especially water and weather resistance with a water-based system. This is because a water based system would have a higher propensity to degrade or be washed away with time. It is for this reason that the joint fillers for exterior sheathing boards have largely been non-water based materials such as silyl-terminated polyethers. Where water based formulations have been used, they have had to be very high in latex content and correspondingly low in water content. Formulations that are high in organic polymer content (such as STPE and water based acrylic latexes) invariably come with handling difficulties such as poor washability and undesirable odors. This invention addresses the above problems and relates to an exterior joint compound that has a significantly lower acrylic latex content, yet displays exceptional water, weather, and mold resistance. This inventive joint compound contains a specially formulated wax emulsion designed to impart exceptional water repellency to coatings formulations. Particularly, the invention relates to exterior joint compounds that are water-based but with water repellant and anti-microbial character. Because the exterior of a building especially during construction is exposed to sever weather conditions such as rain and moisture, it is counterintuitive to envision using a water-based system that can handle the harshness of the exterior environment especially when such joints serve as the weak link for moisture entry. Wax emulsions have been used in composite wallboard (e.g., gypsum wallboard) for many years for the interior in the building industry. For example, wax emulsions sold under the trade name AQUALITE® by Henry Company, and several wax emulsion formulations are disclosed in the prior art, such as U.S. Pat. No. 5,437,722. However, wax emulsions have not been used in exterior joint compounds. The inventors of this invention have arrived at a water based formulation that includes wax-based emulsion systems as one component to be used a joint compound between the exterior sheathing wallboard panels."}
{"text": "The present invention relates to an automatic water supply device for animal drinking and for sprinkling in gardening or horticulture, and more particularly, to an improved automatic water supply device which utilizes siphon action and which supplies water at long-timed intervals in response to lowering of water level in the water tray down to a predetermined level.\nIn most of the prior art automatic water supply devices or systems for animal drinking, a water supply control means such as a valve, ball tap or the like are usually arranged within or adjacent to a water feeding tray or vessel so that said valve, etc., is actuated when it is touched and moved by the mouth or nose of an animal wishing to drink water.\nTherefore, malfunction of this kind of device frequently occurs due to unnecessary engagement of the animal's mouth or nose with such a water supply control means by mere mischief, etc., thereby causing overflow of water on the floor of animal's enclosure without being effectively consumed by the animals, and further making the floor wet and dirty.\nFurther, with respect to the conventional water supply devices, the water control device such as a valve, etc., is soon damaged by corrosion, and the water tray becomes undesirably large in order to receive the control device therein.\nIn view of the abovementioned drawbacks, I have proposed an automatic water supply device utilizes a siphon means, in Japanese Patent Application No. 37062/77 filed on Mar. 31, 1977 (Japanese Laying Open Publication No. 122570/78) and in the corresponding U.S. Pat. application Ser. No. 891,172 filed on Mar. 29, 1978, now U.S. Pat. No. 4,202,294, issued May 13, 1980.\nThe inventor's abovementioned prior invention has solved most of the drawbacks of the conventional devices. However, it has been found that some matters should still be improved with respect to its siphon action.\nIn the device of the inventor's prior invention, if the size and the arrangement of the inversely J-shaped siphon tube are not proper, water droplets remain in the passage of the siphon tube after the end of the siphoning operation in a manner such that some of the water droplets block an intermediate portion of the passage of the siphon tube, while additional droplets are suspended in the inlet tube portion extending from the top portion of the inversely J-shape siphon tube toward its open inlet end. If this happens, when water is poured into the flow control container for the next siphoning supply of water, the water level in the inlet tube portion of the siphon tube may rise in advance of the water level outside of the siphon tube, and due to this discrepancy between the water levels inside and outside of the siphon tube, the siphoning operation may be inadvertently triggered at an early stage before the container has been filled with a predetermined amount of water.\nThis sucking up of water is likely to be accompanied by a small amount of dregs of animal food. The water and dregs are raised, gradually by repetition, through the water feed pipe (the long leg of the siphon) until they finally enter into the flow control container, thereby producing dirty water in the flow control container. Although this contamination is sometimes negligible, it is still desirable that it should be definitely avoided."}
{"text": "This invention relates generally to a ceramic composition and to a method for producing said composition.\nConventional ceramic processing typically comprises ball milling a slurry, which contains a ceramic material and a solvent, to reduce the particle size of the ceramic material. The milled slurry is dried, such as by spray drying or pan drying, and the dried ceramic composition is then formed into articles.\nIn the fabrication of ceramic articles, a binder may be employed. A good binder must satisfy several criteria. It should be soluble in the slurry containing the ceramic and it should have a low solution viscosity thereby permitting high solids content in the slurry and facilitating ball milling. Also, a good binder should be compatible with other additives and processing steps. Since the binder is not a desirable component of a final ceramic article, it should be easily removable at some point in the processing sequence. In addition, the binder should impart green strength to the green fabricated article."}
{"text": "The present invention relates to tracheostoma valves. Tracheostoma valves are currently used in conjunction with speech prosthesis devices and with speaking or fenestrated tracheostomy tubes as a means of diverting exhaled air which would otherwise pass through the tracheostoma to enable the user to speak. It is envisioned the valves may be used as rehabilitation devices for patients seeking to refine esophageal speech and as part of other surgical prosthetic forms of speech rehabilitation.\nIt is a primary object of the present invention to provide a novel tracheostoma valve.\nAnother object of the present invention is to provide a novel tracheostoma valve of the disc-valve type which is constituted and arranged in a novel and expeditious manner.\nTracheostoma valves having a flexible valve diaphragm have been heretofore known in the art being shown, for example, in U.S. Pat. No. 4,325,366, issued Apr. 20, 1982 to Carl J. Tabor, and in an application for U.S. Pat., Ser. No. 316,055, filed Oct. 29, 1981, by Eric D. Blom and Mark I. Singer. However, such prior art valve structures possess several inherent disadvantages, such as, for example, experiencing difficulties in remaining open to accommodate a wide range of pressure and flow conditions associated with breathing and physical exercising, and having flexible valving diaphragm members which are extremely difficult to produce in mass production with desired precision. Also, such flexible valve diaphragms include a crease therein when in the open position, which crease does not provide an adequate seal when the diaphragm is then positioned in the closed position. Also, such valve structures readily fatigue in that they change over time and do not retain uniform physical properties during usage, such as, for example, such valve structures become more flexible through usage and thier position at rest varies.\nAn important object of the present invention is to overcome the difficulties of heretofore known tracheostoma valves in a novel and expeditious manner.\nAnother object of the present invention is to provide a novel tracheostoma valve which accommodates normal vegetative breathing during various physical activities, during speaking and during coughing in a novel and expeditious manner.\nA further object of the present invention is to provide a novel tracheostoma valve which has improved sealing characteristics in the closed or speaking position.\nAnother object of the present invention is to provide a novel tracheostoma valve the operating characteristics of which may be readily varied to meet the needs of both the same individual and different individuals using the valve and which permits air volumes and/or airflows to enter the valve and the lungs unobstructed when in the open or inhaling position.\nYet another object of the present invention is to provide a novel tracheostoma valve which may be readily disassembled and assembled for purposes of cleaning, repairing, and the like.\nA further object of the present invention is to afford a novel tracheostoma valve of the aforementioned type which is practical, efficient and reliable in operation and which may be readily and economically produced commercially.\nOther and further objects of the present invention will be apparent from the following description and claims and are illustrated in the accompanying drawings which, by way of illustration, show the preferred embodiment of the present invention and the principles thereof and what we now consider to be the best mode in which we have contemplated applying these principles. Other embodiments of the invention embodying the same or equivalent principles may be used and changes may be made as desired by those skilled in the art without departing from the present invention and the purview of the appended claims."}
{"text": "The present invention is directed toward a method for liquefying target tissue within a body and more particularly, toward a method of heating a biocompatible fluid and presenting such to target tissue, thereby liquefying the target tissue and aspirating it while leaving non-target tissue intact.\nPrior thermal energy sources used in surgery have been xe2x80x9cdry-heatxe2x80x9d devices. Some examples are: thermal lasers, ultrasound, microwave, radio frequency, and electrosurgery devices. These dry-heat energy sources transfer heat from a surgical instrument to a tissue primarily by the heat transfer modalities of conduction and radiation. Thermal lasers, ultrasound, and microwave energy sources can liquefy tissue, but only at very high energy levels. If any of these dry-heat energy sources are placed in direct contact with target surgical tissue, there is no tissue differentiation. That is, healthy as well as unhealthy tissue will be liquefied. For example, collapsing cavitation bubbles in front of an ultrasound surgical instrument release a thermal energy of 13,000xc2x0 F. per bubble. With such a high thermal energy input there is no differentiation of tissue within the surgical tissue field. That is, all tissue target and non-target, within the target field melts or liquefies. The thermal energy in thermal lasers and ultrasound is inherently high, so the heat cannot be xe2x80x9cturned downxe2x80x9d or reduced. In microwave, radio frequency, and electrosurgery the level of thermal energy may be reduced but in so doing, liquefaction of tissue is not achieved at all. Rather, cauterization or coagulation of the tissue results.\nU.S. Pat. No. 4,924,863 discloses a method for removing artherosclerotic plaque from a patient by heating the plaque under pressure for a limited period of time. When the plaque is heated, it liquefies without causing immediate death of the underlying tissue. The liquefied plaque is then aspirated from the body through a catheter. The means used for heating the plaque is microwave or radio frequency energy. As discussed above, however, in an effort to control the level of energy being produced, all of the plaque may not liquefy.\nU.S. Pat. No. 5,540,679 discloses the use of a balloon catheter for heating tissue in a patient\"\"s body. A heating device is located within the balloon and is arranged for heating fluid inside the balloon. When the fluid is heated, unwanted tissue, such as a tumor or an enlarged prostate gland, is heated by thermal conduction from the fluid through a wall of the balloon. The heat destroys the tissue which is eventually absorbed into the patient\"\"s body. This method uses conduction anand as explained above, such a method creates such a high level of energy that there is no tissue differentiation. As a result, healthy tissue may be heated and ultimately destroyed.\nU.S. Pat. No. 4,886,491 discloses a method of liposuction using an ultrasonic probe. An ultrasonic probe tip is vibrated at a high frequency and a low amplitude. This method separates the fatty tissue and creates heat which melts some of the fatty tissue. A saline irrigating solution is applied to the area which emulsifies the melted fatty tissue. The emulsified solution is then aspirated. Again, because of the high thermal energy input into the system there is no differentiation of tissue into target and non-target tissue.\nThe present invention is designed to overcome the deficiencies of the prior art discussed above. It is an object of this invention to provide a method for liquefying target tissue within a living body while leaving non-target tissue intact, and aspirating the liquefied tissue as it is liquefied.\nIt is another object of the invention to provide a method for liquefying target tissue by heating a biocompatible fluid and presenting the same to the target tissue.\nIt is a further object of the invention to provide a method for liquefying fatty tissue and aspirating the same from a patient.\nIn accordance with the illustrative embodiments, demonstrating features and advantages of the present invention, there is provided a method for liquefying target tissue within a human body. The method includes heating a biocompatible fluid, presenting this fluid to the target tissue which liquefies the tissue, while leaving non-target tissue intact, and aspirating the liquefied tissue as soon as it is liquefied.\nThe present invention discloses a novel energy source which allows for efficient, safe, minimally invasive, and cost-effective surgery for surgical procedures where the goal of the surgery is to remove unwanted tissue from the body, within a surgical tissue field. The present invention also allows for the differentiation of target tissue from non-target tissue. In other words, the surgical tissue field is defined as the actual anatomical landscape with which the surgeon is physically interacting. Surgical tissue fields are comprised of multiple homogeneous tissue groups. Each of these homogeneous tissue groups has its own melting point. With the present invention, the melting point of a particular tissue group or unwanted tissue is determined in order to liquefy the tissue and rid the patient of the same."}
{"text": "1. Field of the Invention\nThe present invention relates to a tool suspension device having receiving element which is slidable.\n2. Description of the Prior Art\nA conventional tool suspension device usually includes anti-theft mechanisms. For example, a slidable receiving element is included. The receiving element is prohibited from being removed by a hindering mechanism blocking the receiving element. As disclosed in patent TW M320343, the tool suspension device includes a main body, a sliding track, and a plurality of receiving elements. The receiving elements are slidably disposed on the sliding track. Two ends of the sliding track are positioned to the main body by covers. The receiving elements are unable to be removed from the two ends of the sliding track freely.\nHowever, the device has a complicated structure to have too much components. More specifically, the main body, the sliding track, the receiving elements, and at least two covers are necessary in this disclosure. Furthermore, bolts are also necessary for fixing. If any one of the covers is lost, the receiving element is unable to be blocked anymore. On the other hand, if any one of the covers is removed by tools, the receiving element is able to be easily removed. Besides, due to the specific positions of the covers, the length of the sliding track is constant."}
{"text": "The invention relates to apparatus and methods for carrying out electrodeionization to purify water.\nElectrodeionization is a process for removing ions from liquids by sorption of these ions into a solid material capable of exchanging these ions for either hydrogen ions (for cations) or hydroxide ions (for anions) and simultaneous or later removal of the sorbed ions into adjacent compartments by the application of an electric field. (See Glueckauf, E., \"Electro-Deionization Through a Packed Bed\", December 1959, pp. 646-651, British Chemical Engineering for a background discussion.) The hydrogen and hydroxide ions needed to drive the ion exchange process are created by splitting of water molecules at the interface of anion and cation exchanging solids which contact each other in the orientation that depletes the contact zone of ions, when in the presence of an electric field. This orientation requires that the anion exchanging material face the anode and the cation exchanging material face the cathode. The created hydroxide ions enter the anion exchanging material, and the created hydrogen ions enter the cation exchanging material.\nThe electrodeionization process is commonly carried out in an apparatus consisting of alternating diluting compartments and concentrating compartments separated by anion permeable and cation permeable membranes. The diluting compartments are filled with porous ion exchanging solid materials through which the water to be deionized flows. The ion exchanging materials are commonly mixtures of cation exchanging resins and anion exchanging resins (e.g., U.S. Pat. No. 4,632,745), but alternating layers of these resins have also been described (e.g., U.S. Pat. Nos. 5,858,191 and 5,308,467). Ion exchanging materials consisting of woven and non-woven fibers have also been described. E.g., U.S. Pat. No. 5,308,467 describes a fabric in which bundles of cation-exchange fibers and are woven alternately with bundles of anion-exchange fibers, and U.S. Pat. No. 5,512,173 describes a cloth containing cation exchange fibers, anion exchange fibers and ionically inactive fibers. The compartments adjoining the diluting compartment into which the ions are moved by the applied electric field, called concentrating compartments, may be filled with ion exchanging materials or with inert liquid permeable materials. An assembly of one or more pairs of diluting and concentrating compartments, referred to as a \"cell pair\", is bounded on either side by an anode and a cathode which apply an electric field perpendicular to the general direction of liquid flow.\nThe diluting compartments are each bounded on the anode side by an anion permeable membrane and on the cathode side by a cation permeable membrane. The adjacent concentrating compartments are each correspondingly bounded by a cation permeable membrane on the anode side and an anion permeable membrane on the cathode side. The applied electric field causes anions to move from the diluting compartment across the anion permeable membrane into the concentrating compartment nearer the anode and cations to move from the diluting compartment across the cation permeable membrane into the concentrating compartment nearer the cathode. The anions and cations become trapped in the concentrating compartments because the movement of anions toward the anode is blocked by a cation permeable membrane, and the movement of cations toward the cathode is blocked by an anion permeable membrane. A flow of water is set up to remove the ions from the concentrating compartments. The net result of the process is the removal of ions from the water stream flowing through the diluting compartments and their concentration in the water flowing through the concentrating compartments.\nThe removal of the ions from the diluting compartment is a multi-step process involving diffusive steps as well as electrically driven steps. First, it is clear that the movement of ions directly from the diluting solution across the bounding membranes, under the influence of the applied electric field, contributes insignificantly to the overall removal of these ions, because the concentration of ions in the diluting solution is typically 1,000 to 100,000 times smaller than the concentration of ions in the solid ion exchanging materials. While the mobility of ions in the solid material may be on the order of 20 times smaller than their mobility in the solution, the electric field acting on the ions in the two phases is the same, so the product of mobility times concentration times electric field strength, which determines the rate of ion removal, is 50 to 5,000 times as large in the solid ion exchanging material.\nGlueckauf showed that the mechanism of ion removal from the diluting compartment solution includes two steps. The first step is the diffusion of cations to the cation exchanging solids and the diffusion of anions to the anion exchanging solids. The second step is electrical conduction within the solids phases to the bounding membranes of the diluting compartment. Because the concentration of ions in ion exchanging solids is so high, the process that controls the overall removal of ions is their rate of diffusion from the solution to the surface of the ion exchanging solids. This diffusion rate is a function of three factors; the diffusion rate is proportional to surface area between the ion exchanging solids and the flowing solution, inversely proportional to the thickness of the liquid layer through which the ions must diffuse, and proportional to the difference in concentration of the ions in the bulk of the diluting solution and their concentration next to the ion exchanging solid. In order to achieve high rates of ion removal, the product of the above three factors should thus be as high as possible. The ratio of the surface area to the diffusion distance is inversely proportional to the characteristic dimension of the ion exchanging solid material; the characteristic dimension is particle radius for ion exchange resins and is fiber diameter for ion exchange fibers. In designing electrodeionization apparatus, this characteristic dimension can be made as small as possible, commensurate with avoidance of excessive pressure drops or plugging by particles in the water to be treated. Particle diameters on the order of 500 to 600 micrometers are typical, and fiber diameters can be on the order of several tens of microns.\nAs noted above, the third factor controlling the rate of ion removal is the difference in concentration of the ion being removed between the bulk of the solution and its concentration in the liquid adjacent to the surface of the ion exchanging solid where it is being exchanged for either a hydrogen or a hydroxide ion. The concentration of the ion in question at the surface of the ion exchanging solid is in equilibrium with the concentration of that ion in the solid. For cations, the equilibrium concentration is approximately equal to the ratio of the cation concentration to the hydrogen ion concentration in the cation exchanging solid times the concentration of the cation in solution. For anions, the equilibrium concentration is approximately equal to the ratio of the anion concentration to the hydroxide ion concentration in the anion exchanging solid times the concentration of the anion in solution. In order for this equilibrium concentration to be low, and the rate controlling concentration difference to be large, the cation exchanging solid should be predominantly in the hydrogen form, and the anion exchanging solid should be predominantly in the hydroxide form. In fact, if the two solids are completely in the ionic form rather than in the hydrogen or hydroxide form, there is no concentration difference, and ions will not be removed by this diffusive mechanism.\nIn order for the ion exchanging solids to be predominantly in the hydrogen and hydroxide forms, the so-called \"regenerated forms,\" the rate of hydrogen ion and hydroxide ion creation (water splitting) must be both high and spatially uniform. A high average rate of water splitting can be achieved by applying a high voltage drop across the diluting compartment. With equinormal mixtures of ion exchange particles, voltages of between 1 and 5 volts are adequate for the purpose. The achievement of a uniform distribution of water splitting is a more difficult problem and much effort has gone into designing structures that achieve this (e.g., U.S. Pat. Nos. 5,858,191, 5,868,915 and 5,308,467). The random nature of mixtures of cation and anion exchanging particles tends to cause some portion of the particles to be regenerated to a needlessly high degree and others inadequately regenerated. Water flowing through the regions of inadequately regenerated material will be inadequately purified. The essence of the difficulty that existing approaches have had in dealing with the problem is that the number of contacts between cation exchanging and anion exchanging material where water splitting can take place is limited by the relatively large characteristic dimension of the ion exchanging material. This results in regions of inadequately regenerated resin between the water splitting sites."}
{"text": "Vehicle windows (e.g., windshields, backlites, sunroofs, and sidelites) are known in the art. For purposes of example, vehicle windshields typically include a pair of bent glass substrates laminated together via a polymer interlayer such as polyvinyl butyral (PVB). It is known that one of the two glass substrates may have a coating (e.g., low-E coating) thereon for solar control purposes such as reflecting R and/or UV radiation, so that the vehicle interior can be more comfortable in certain weather conditions. Conventional vehicle windshields are made as follows. First and second flat glass substrates are provided, one of them optionally having a low-E coating sputtered thereon. The pair of glass substrates are washed and booked together (i.e., stacked on one another), and then while booked are heat bent together into the desired windshield shape at a high temperature(s) (e.g., 8 minutes at about 600-625 degrees C.). The two bent glass substrates are then laminated together via the polymer interlayer to form the vehicle windshield.\nInsulating glass (IG) window units are also known in the art. Conventional IG window units include at least first and second glass substrates (one of which may have a solar control coating on an interior surface thereof) that are coupled to one another via at least one seal(s) or spacer(s). The resulting space or gap between the glass substrates may or may not be filled with gas and/or evacuated to a low pressure in different instances. However, many IG units are required to be tempered. Thermal tempering of the glass substrates for such IG units typically requires heating the glass substrates to temperature(s) of at least about 600 degrees C. for a sufficient period of time to enable thermal tempering.\nOther types of coated articles also require heat treatment (HT) (e.g., tempering, heat bending, and/or heat strengthening) in certain applications. For example and without limitation, glass shower doors, glass table tops, and the like require HT in certain instances.\nDiamond-like carbon (DLC) is sometimes known for its scratch resistant properties. For example, different types of DLC are discussed in the following U.S. Pat. Nos. 6,303,226; 6,303,225; 6,261;693; 6,338,901; 6,312,808; 6,280,834; 6,284,377; 6,335,086; 5,858,477; 5,635,245; 5,888,593; 5,135,808; 5,900,342; and 5,470,661, all of which are hereby incorporated herein by reference.\nIt would sometimes be desirable to provide a window unit or other glass article with a protective coating including DLC in order to protect it from scratches and the like. Unfortunately, DLC tends to oxidize and burn off at temperatures of from approximately 380 to 400 degrees C. or higher, as the heat treatment is typically conducted in an atmosphere including oxygen. Thus, it will be appreciated that DLC as a protective overcoat cannot withstand heat treatments (HT) at the extremely high temperatures described above which are often required in the manufacture of vehicle windows, IG window units, glass table tops, and/or the like. Accordingly, DLC cannot be used alone as a coating to be heat treated, because it will oxidize during the heat treatment and substantially disappear as a result of the same (i.e., it will burn off).\nCertain other types of scratch resistant materials also are not capable of withstanding heat treatment sufficient for tempering, heat strengthening and/or bending of an underlying glass substrate.\nAccordingly, those skilled in the art will appreciate that a need in the art exists for a method of making a scratch resistant coated article that is capable of being heat treated (HT) so that after heat treatment the coated article is still scratch resistant. A need for corresponding coated articles, both heat treated and pre-HT, also exists."}
{"text": "1. Field of the Invention\nThe present invention relates to an image generating apparatus for generating an image on recording paper through an electrophotographic or xerographic process.\n2. Description of the Related Art\nAs the processing speed of personal computers and workstations has increased on recent years, it has also been increasing demand to speed up the processing of an image generating apparatus based on a xerographic process. With such an image generating apparatus as mentioned above, when a polygon mirror included in an exposure optical system is increased in its rotational speed, it is possible to speed up its printing operation. In the prior art, however, the increased rotational speed of the polygon mirror has already reached its limit. For this reason, as another additional speeding-up measure, a plurality of semiconductor lasers for forming an electrostatic latent image on a photoconductor drum are provided to realize a simultaneous recording system.\nAn arrangement and operation of a prior art image generating apparatus having a plurality of semiconductor lasers will be explained in the following. FIG. 11 is an arrangement of an exposure optical system having a plurality of semiconductor lasers therein in the prior art. The exposure optical system of FIG. 11 includes first and second semiconductor lasers 21 and 22, which are positioned so that laser beams emitted from the lasers 21 and 22 are crossed perpendicularly to each other. A collimator lens 23 for making the input light beams parallel is positioned in a beam exit direction of the first semiconductor laser 21, and a collimator lens 24 is similarly positioned on a beam exit direction of the second semiconductor laser 22. Disposed on light exit sides of the collimator lenses 23 and 24 is a beam splitter 25 which functions to deflect the laser beam of the first semiconductor laser 21 to a direction perpendicular thereto, transmit the laser beam of the second semiconductor laser 22 therethrough and guide the beams onto a surface of a polygon mirror 26. The polygon mirror 26, which is disposed on a rotatary axis of a polygon motor (not shown in the drawing) rotating at a high speed, acts to cause the laser beams received from the first and second semiconductor lasers 21 and 22 to be scanned on a photoconductor drum 1. Disposed between the polygon mirror 26 and the photoconductor drum 1 are on f lens 27 for reducing the diameter of the laser beams reflected on the polygon mirror 26 to a predetermined value as well as on a reflecting mirror 28.\nWith this image generating apparatus, image data corresponding to 2 lines can be recorded at a time in a main scanning direction B of the photoconductor drum 1 from the 2 laser beams received from the first and second semiconductor lasers 21 and 22.\nFurther, FIG. 12 shows a control block diagram of a prior art image generating apparatus having 2 semiconductor lasers. FIG. 13 is a detailed circuit diagram of a first image data selective output means in the prior art image generating apparatus. In FIG. 12, an image data generating means 31 generates image data (bit map data) on the basis of image information received from a host computer (not shown). A memory means 32 stores therein the image data generated by the image data generating means 31. Further, a first image data selective output means 37 and a second image data selective output means 38 act to extract the image data from the memory means 32 and output it to a first exposure energy control means 35 and a second exposure energy control means 36, respectively. The first and second exposure energy control means 35 and 36 respectively control exposure energy (light emission time, light emission power) of the first and second semiconductor lasers 21 and 22 incorporated in first and second exposure means 33 and 34 (which will be explained later).\nThe first and second exposure means 33 and 34, which include the first and second semiconductor lasers 21 and 22, respectively, irradiate laser beams on the photoconductor drum 1 to form an electrostatic latent image on the photoconductor drum 1.\nAn output control means 39, on the basis of a printing operation reference signal, controls the operation of the sequential output of the image data stored in the memory means 32 to the first and second image data selective output means 37 and 38. Further, a clock generating means 40 generates a clock (which will be referred to as the video clock) that is used as a reference to the output operation of the first and second semiconductor lasers 21 and 22.\nExplanation will now be made as to the operation of the image generating apparatus having the above arrangement. In FIG. 11, image data is optically modulated by the first and second semiconductor lasers 21 and 22 so that output laser beams of the respective semiconductor lasers are converted by the collimator lenses 23 and 24 to collimated or parallel light and then input to the beam splitter 25. The beam splitter 25 in turn deflects by 90 degrees the laser beam received from the first semiconductor laser 21 and then guides it to the polygon mirror 26; whereas the beam splitter 25 transmits the laser beam received from the second semiconductor laser 22 therethrough and then guides it to the polygon mirror 26. The polygon mirror 26, while rotated by the polygon motor (not shown), scans the laser beams received from the first and second semiconductor lasers 21 and 22 at a predetermined angle. The laser beams are further input to the f lens 27 where the laser beams are reduced in diameter to a predetermined diametered single beam, and then the single beam is scanned by the reflecting mirror 28 on the photoconductor drum 1 in the direction (main scanning direction) shown by an arrow B. In this case, the photoconductor drum 1 is rotating at a constant speed Vp (mm/sec.) in a direction shown by an arrow A.\nFIG. 14 is a diagram for explaining how image data is generated in the prior art image generating apparatus. As shown in FIG. 14, picture elements (pixels) 30 each having an illustrated diameter are formed by the laser beam issued from the second semiconductor laser 22 on the first raster scan line, while picture elements (pixels) 29 each having the same diameter as that of the pixel 30 are formed by the laser beam issued from the first semiconductor laser 21 on the second raster scan line. The above operation is sequentially repeated in a direction shown by an arrow C so that the laser beam issued from the second semiconductor laser 22 forms pixels of the odd-numbered raster scan lines, whereas the laser beam issued from the first semiconductor laser 21 forms pixels of the even-numbered raster scan lines. As a result, printing can be realized at a speed of twice as fast as the case of using a single semiconductor laser.\nExplanation will next be made as to specific control operation when image data having a main scan (horizontal scanning) resolution of 600 dpi and a feed scan (vertical scanning) or sub-scan resolution of 600 dpi is output to the photoconductor drum 1. FIG. 15 is a timing chart of the first image data selective output means of the prior art image generating apparatus, and FIG. 16 is a diagram showing a relationship between pixel formation and laser emission time in the prior art image generating apparatus.\nIn FIG. 12, the image data generating means 31, on the basis of image information received from the host computer (not shown), generates image data (bit map data) and stores it in the memory means 32. On the basis of a printing operation reference signal (not shown), the output control means 39 controls the memory means 32 to cause the memory means 32 to output the first and second raster scan lines of the image data therefrom to the second and first image data selective output means 38 and 37 respectively. That is, as shown in FIG. 13, a semiconductor memory (DRAM) 41 within the memory means 32 is connected to a parallel/serial converter 42 within the first image data selective output means 37 through a bus 8 bits of D7-D0. The output control means 39 generates a predetermined address and sends it to the DRAM 41 to read out image data at the address, and loads the read-out image data to the parallel/serial converter 42 under a signal LD.\nThe image data of D7-D0 loaded into the parallel-to-serial converter 42, as shown in FIG. 15, in synchronism with a video clock f received from the clock generating means 40, serially outputs the image data (starting from D7 and ending in D0) from a serial output terminal Q7 to the first exposure energy control means 35 as a laser beam emission signal of the second raster scan line.\nFurther, the output control means 39 outputs the next address to read out image data at the address from the DRAM 41. The similar operation to the above is repeated to cause the parallel-to-serial converter 42 to sequentially output the image data of the fourth raster scan line as a laser beam emission signal.\nSimilarly, the second image data selective output means 38 sequentially outputs the image data of the first raster scan line to the second exposure energy control means 36 as serial image data.\nThe first and second exposure energy control means 35 and 36, according to the received image data, send the respective laser beam emission signals of predetermined emission time and power to the first semiconductor laser 21 of the first exposure means 33 and also to the second semiconductor laser 22 of the second exposure means 34 for modulation, respectively. As shown in FIG. 15, the laser beam emission time of the first and second semiconductor lasers 21 and 22 is usually set at a period time of the video clock f, so that such a predetermined size (beam diameter) of pixel dots as shown in FIG. 16 is formed. In the drawing, the black dot or circle denotes the presence of the pixel output and the white dot or circle denotes the absence of the pixel output.\nIn the case of the prior art image generating apparatus having the aforementioned arrangement, in order to obtain a high quality of such text image as characters or lines, it is necessary to increase its pixel density or resolution. FIG. 17 is a timing chart of the first image data selective output means of the prior art image generating apparatus when the apparatus is operated in a high-density record mode, and FIG. 18 shows a relationship between the laser beam emission signal and formed pixels with the same high resolution. For example, when the resolution in the main scanning direction B is changed from 600 dpi to 1200 dpi as shown in FIG. 18, it is necessary to set the frequency of the video clock at twice (2 fHz) as shown in FIG. 17.\nHowever, this also requires the other circuits to operate at a speed of twice the previous operational speed, which also increases noise such as unnecessary irradiation. In this way, this involves a problem in which the circuit costs become high.\nIn order to obtain a high quality of image such as a photograph, further, it is necessary to increase or enhance its pixel density or gray shade property (gradient). For enhancing the gray shade property, it is required to divide exposure energy into a plurality of stages of energy. However, the division of exposure energy into such a plurality of stages by one scan of the semiconductor laser involves the complication of the control circuit, thus disadvantageously increasing the circuit costs.\nIt is accordingly an object of the present invention to provide an image generating apparatus which can produce a high quality of image without the need for increasing the frequency of a video clock and can provide gray shade display with use of a simple control circuit.\nIn accordance with the present invention, there is provided an image generating apparatus which has a plurality of exposure means for optically modulating image data to record a plurality of main scanning lines of the image data on a recording medium through a single main scanning operation, and which comprises image data selective output means for selecting the image data corresponding to the plurality of exposure means from the image data forming one of the main scanning lines and outputting the selected image data to the respective exposure means; and control means for controlling the operation of the image data selective output means in such a manner that the plurality of exposure means perform sequential scanning operation over the one main scanning line based on the image data issued from the image data selective output means to record the image data.\nIn the image generating apparatus of the present invention, one of the main scanning lines of image data is recorded through a plurality of scanning operations of the plurality of exposure means. For this reason, even when the number of pixels of one main scanning line is increased, the number of image data to be recorded for one main scanning line by one exposure means can be maintained as in the prior art. Accordingly, high-density image formation can be implemented through a plurality of recording operations of the plurality of exposure means over the image data without the need for increasing the frequency of a video clock for generating the timing of irradiation or illumination of an exposure light beam from one exposure means.\nIn accordance with the present invention, there is also provided an image generating apparatus which has a plurality of exposure means for optically modulating image data to record a plurality of main scanning lines of the image data on a recording medium through a single main scanning operation, and which comprises image data division/output means for dividing the image data forming one of the main scanning lines into predetermined gray shade levels of image data and outputting the divided image data to the respective exposure means; and control means for controlling the operation of the image data division/output means in such a manner that the plurality of exposure means perform sequential scanning operation over the one main scanning line based on the image data issued from the image data divisional/output means to record the image data.\nIn the image generating apparatus of the present invention, the image data division/output means divides the image data into different gray shade levels of image data and outputs them to the plurality of exposure means. The respective exposure means form pixels allowing the gray shade display according to the gray shade levels of the received image data. In addition, image data are recorded based on the different-level gray shade display through a plurality of scanning operations of the plurality of exposure means over one main scanning line. Since one exposure means can perform the gray shade display as controlled to a constant exposure condition, the gray shade display can be realized through easy-to-control operation.\nIn accordance with a first aspect of the present invention, there is provided an image generating apparatus which has a plurality of exposure means for optically modulating image data to record a plurality of main scanning lines of the image data on a recording medium through a single main scanning operation, and which comprises image data selective output means for selecting the image data corresponding to the plurality of exposure means from the image data forming one of the main scanning lines and outputting the selected image data to the respective exposure means; and control means for controlling operation of the image data selective output means in such a manner that the plurality of exposure means perform sequential scanning operation over the one main scanning line based on the image data issued from the image data selective output means to record the image data.\nThus, high-density image formation can be carried out without the need for increasing the frequency of a video clock signal for controlling the timing of exposure operations of the respective exposure means.\nIn accordance with a second aspect of the present invention, in the image generating apparatus of the first aspect, a plurality of the image data selective output means are provided for the plurality of exposure means, and the control means performs the control operation in such a manner that as a main scanning position of the plurality of exposure means advances by one line in a sub-scanning direction perpendicular to the main scanning direction, the plurality of exposure means simultaneously record the image data on a plurality of main scanning lines based on the image data issued from the image data selective output means.\nTherefore, when the plurality of exposure means are advanced by an amount of each one line in the sub-scanning direction through the repetitive main scanning operation, one main scanning line of image data can be recorded through the scanning operations of the plurality of exposure means, and thus high-density image formation can be realized without the need for increasing the frequency of the video clock signal.\nIn accordance with a third aspect of the present invention, in the image generating apparatus of the first or second aspect, the image data selective output means includes change-over means for switching between a first mode in which a plurality of main scanning lines of image data to be simultaneously recorded are output to the associated exposure means and a second mode in which the image data selective output means selects image data corresponding to the plurality of exposure means from the image data of one main scanning line and outputs the selected image data to the respective exposure means.\nThus, selection can be made between the normal operation of image formation of plural lines and the high-density image formation operation.\nIn accordance with a fourth aspect of the present invention, there is provided an image generating apparatus which has a plurality of exposure means for optically modulating image data to record a plurality of main scanning lines of the image data on a recording medium through a single main scanning operation, and which comprises image data division/output means for dividing the image data forming one of the main scanning lines into predetermined gray shade levels of image data and outputting the divided image data to the respective exposure means; and control means for controlling operation of the image data division/output means in such a manner that the plurality of exposure means perform sequential scanning operation over the one main scanning line based on the image data issued from the image data division/output means to record the image data.\nTherefore, the gray shade display using the plurality of exposure means can be realized through simple control operation while eliminating the need for changing the exposure conditions of one exposure means.\nIn accordance with a fifth aspect of the present invention, in the image generating apparatus of the fourth aspect, the plurality of exposure means, on the basis of image data received from the image data division/output means, irradiate exposure light beams having different exposure energies on the recording medium and overlap the exposure light beams issued from the plurality of exposure means with respect to predetermined pixels to record the image data of the different gray shade levels.\nThus, image gray shade display can be realized by making different the exposure energies of the plurality of exposure means.\nIn accordance with a sixth aspect of the present invention, in the image generating apparatus of the fourth or fifth aspect, a plurality of image data division/output means are provided for the plurality of exposure means.\nTherefore, when the respective image data division/output means output image data to the associated exposure means according to the gray shade levels, gray shade display can be realized.\nIn accordance with a seventh aspect of the present invention, in the image generating apparatus of any of the fourth to sixth aspects, the image data division/output means includes change-over means for switching between a first mode in which a plurality of main scanning lines of image data to be simultaneously recorded are output to the associated exposure means and a second mode in which the image data division/output means divides image data of one of the main scanning lines into predetermined gray shade levels of image data and outputs the divided image data to the respective exposure means.\nThus, change-over can be carried out between the normal mode of image formation by the simultaneous scanning operation of a plurality of lines and the image formation mode based on the gray shade display."}
{"text": "1. Field of Invention\nThe field of the currently claimed embodiments of this invention relates to optical coherence tomography systems, and more particularly to motion-compensated optical coherence tomography systems.\n2. Discussion of Related Art\nOptical coherence tomography (OCT) has been viewed as an “optical analogy” of ultrasound sonogram (US) imaging since its invention in early 1990's (D. Huang, E. A. Swanson, C. P. Lin, J. S. Schuman, W. G. Stinson, W. Chang, M. R. Hee, T. Flotte, K. Gregory, C. A. Puliafito, and J. G. Fujimoto, “Optical coherence tomography,” Science, vol. 254, pp. 1178-1181, 1991). Compared to the conventional image-guided interventions (IGI) using modalities such as magnetic resonance imaging (MRI), X-ray computed tomography (CT) and ultrasound (US) (T. Peters and K. Cleary, Image-Guided Interventions: Technology and Applications, Springer, 2008), OCT has much higher spatial resolution and therefore possesses great potential for applications in a wide range of microsurgeries, such as vitreo-retinal surgery, neurological surgery and otolaryngologic surgery.\nAs early as the late 1990's, interventional OCT for surgical guidance using time domain OCT (TD-OCT) at a slow imaging speed of hundreds of A-scans/s has been demonstrated (S. A. Boppart, B. E. Bouma, C. Pitris, G. J. Tearney, J. F. Southern, M. E. Brezinski, J. G. Fujimoto, “Intraoperative assessment of microsurgery with three-dimensional optical coherence tomography,” Radiology, vol. 208, pp. 81-86, 1998). Thanks to the technological breakthroughs in Fourier domain OCT (FD-OCT) during the last decade, ultrahigh-speed OCT is now available at >100,000 A-scan/s. For example, see the following:                B. Potsaid, I. Gorczynska, V. J. Srinivasan, Y. Chen, J. Jiang, A. Cable, and J. G. Fujimoto, “Ultrahigh speed Spectral/Fourier domain OCT ophthalmic imaging at 70,000 to 312,500 axial scans per second,” Opt. Express, vol. 16, pp. 15149-15169, 2008.        R. Huber, D. C. Adler, and J. G. Fujimoto, “Buffered Fourier domain mode locking: unidirectional swept laser sources for optical coherence tomography imaging at 370,000 lines/s,” Opt. Lett., vol. 31, pp. 2975-2977, 2006.        W-Y. Oh, B. J. Vakoc, M. Shishkov, G. J. Tearney, and B. E. Bouma, “>400 kHz repetition rate wavelength-swept laser and application to high-speed optical frequency domain imaging,” Opt. Lett., vol. 35, pp. 2919-2921, 2010.        B. Potsaid, B. Baumann, D. Huang, S. Barry, A. E. Cable, J. S. Schuman, J. S. Duker, and J. G. Fujimoto, “Ultrahigh speed 1050 nm swept source/Fourier domain OCT retinal and anterior segment imaging at 100,000 to 400,000 axial scans per second,” Opt. Express, vol. 18, pp. 20029-20048, 2010.        W. Wieser, B. R. Biedermann, T. Klein, C. M. Eigenwillig, and R. Huber, “Multi-Megahertz OCT: High quality 3D imaging at 20 million A-scans and 4.5 GVoxels per second,” Opt. Express, vol. 18, pp. 14685-14704, 2010.        T. Klein, W. Wieser, C. M. Eigenwillig, B. R. Biedermann, and R. Huber, “Megahertz OCT for ultrawide-field retinal imaging with a 1050 nm Fourier domain mode-locked laser,” Opt. Express, vol. 19, pp. 3044-3062, 2011.        \nFor a spectrometer-based SD-OCT, an ultrahigh speed CMOS or CCD line scan camera based system has achieved up to 312,500 line/s in 2008 (Potsaid et al.); while for a swept laser type OCT, >20,000,000 line/s rate was achieved by multi-channel FD-OCT using a Fourier Domain Mode Locking (FDML) laser in 2010 (Wieser et al.).\nFourier-domain optical coherence tomography (FD-OCT) is a high-speed high-resolution three-dimensional imaging modality widely used in biomedical imaging. For OCT to find applications in the interventional imaging area, real-time image processing and display are required.\nOptical coherence tomography (OCT) is a non-invasive, high speed, high-resolution, three-dimensional imaging modality that is widely being used for biomedical application [1, 2]. The real-time non-invasive depth-resolved imaging of tissue structure and flow information provided by OCT can be highly valuable information that can assist physicians in making real-time decisions during surgical procedures such as neurosurgery, tumor resection, microvascular anastomosis, and retinal microsurgery [2-10].\nIn many circumstances, it is more convenient to use a simple hand-held, manually-scanned probe to obtain OCT images of tissues and organs which might otherwise be inaccessible using standard mechanical scanning heads [6]. A hand-held image probe has the following advantages. First, it is small and lightweight, making it easy to access tight spaces. Second, surgeons are intimately familiar with hand-held instruments which can leverage the surgeons' experience and skills with little training. Third, a small hand-held instrument offers greater safety because the surgeon can more easily override or remove the instrument in cases of malfunction [11]. Finally it offers the surgeon great freedom to obtain any image size, for example a larger field-of-view compared to views constrained by apertures of scanning lenses or other endoscopic probes [12-15].\nA hand-held probe, however, poses additional challenges over mechanically-rigid scanners. First, non-uniform motion of the probe during lateral manual scanning will cause image distortion and inaccuracy. Earlier work by Ahmed et al. and more recent work by our lab provide solutions to correct non-uniform scanning speed artifact using de correlation of adjacent A-lines [6, 16]. Second, physiological tremor composed of low and high frequency amplitudes over 100 μm [17] would cause large motion artifacts in acquired images. Third, involuntary motions of a subject may also cause OCT imaging artifacts. Finally, the manual scan across and close to the target surface is highly risky—especially involving fragile tissue. For example, in the context of retinal surgery, the retina is only ˜350 μm for humans, and tearing them can permanently damage eyesight. Scanning while maintaining a larger distance between the probe tip and target surface is not an ideal solution since that degrades image quality and the imaging depth is typically limited to 3-5 millimeters. While motion is in the form of both axial and lateral directions, axial motion is the primary concern due to its direct effects on the image quality. There have been methods used to compensate for the sample surface topology and axial motion in OCT to keep a good system sensitivity and image range, for example the adaptive ranging technique for time-domain OCT (TDOCT) [18] and the reference mirror tracking method for spectral domain OCT (SDOCT) [19]. Common path optical coherence tomography (CP-OCT) is a simple, all-fiber-based technique that shares reference and probe arm which circumvents group velocity dispersion and polarization compensation [20]. It has been widely investigated for various medical applications [21-24]. CP-OCT has also recently been demonstrated for its effectiveness, simplicity, and integrability for surface topology and motion compensation [25] and smart surgical tools [26]. There thus remains a need for improved optical coherence tomography systems that provide motion compensation."}
{"text": "Computing systems often utilize firmware that is stored in a non-volatile memory device, such as a read-only memory (“ROM”) device or a non-volatile random access memory (“NVRAM”) device. The firmware provides program code for performing power-on self tests, booting, and for providing support to the operating system and other functions. When computer systems are designed, there are many different combinations of hardware devices that may be present in the computer system or on a computer system motherboard. To support such a wide variety of hardware combinations, the motherboard manufacturers often provide original equipment manufacturers (“OEMs”) with a reference board and corresponding firmware image that represents the most common motherboard chipset configuration.\nThe reference board provides an OEM with a proof of concept that the OEM will then customize according to their requirements. The OEM may add, remove, or interchange hardware components to and from the reference board according to desired hardware configurations. As various hardware components are added or removed from a computer system motherboard, the firmware image must be updated to reflect the configuration change in order for the modified chipset configuration to properly function. The OEM must then request a firmware image that corresponds to the specific components of the customized motherboard. The customized firmware image is created and provided to the OEM for flashing into a memory device and distribution with the OEM's customized motherboard. This firmware customization process consumes time and resources as the firmware image is created by a manufacturer, delivered to an OEM, returned to the manufacturer for customization, modified by the manufacturer, and returned to the OEM.\nIt is with respect to these considerations and others that the various embodiments of the invention have been made."}
{"text": "1. Field of the Invention\nDuring the Christmas holiday season, tinsel is applied to various forms of Christmas decorations. Tinsel is marketed in packages of elongated strands of tinsel which inevitably become tangled and render the application of individual tinsel strips on articles to be decorated very difficult. In addition, tinsel is also utilized in artistic displays during times other than the Christmas season. Accordingly, a need exists for a structure whereby individual strands of tinsel may be applied to articles to be decorated.\n2. Description of the Prior Art\nVarious forms of ribbon dispensers and wire dispensers have been heretofore provided, but most of these are not suited to dispensing fragile tinsel strands.\nExamples of previously patented wire and ribbon dispensers as well as other structures including some of the general structural and operational features of the instant invention are disclosed in U.S. Pat. Nos. 2,790,925, 3,206,782, 3,382,398, 3,396,888 and 3,811,611."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus, and a method therefor, which is communicatively connected, via a network, to a portable terminal carried by a user, and to a plurality of printing apparatuses provided at various locations, and which outputs, in response to a data print request from the portable terminal, data associated with the data print request to one of the printing apparatuses. More specifically, the invention relates to a data output control apparatus which is suitable for allowing detailed information on a network to be readily obtained.\n2. Description of the Related Art\nIn recent years, information on the Internet is readily available anywhere using, for example, i-mode(trademark) provided by NTT Mobile Communications Network, Inc. (NTT DoCoMo).\nHowever, although such a portable terminal allows information on the Internet to be readily obtained anywhere, because the structure of a display is simplified for purposes, such as reducing the size of, and saving energy for, the portable terminal, the displayed information is rather simplified compared with a display on an ordinary personal computer, which does not satisfy a user with the information service that is provided.\nAccordingly, in order to provide detailed information, while continuing to reduce the size of, and save energy for, the portable terminal, a structure can be provided in which, by combining a portable terminal and a printing apparatus, rough information is displayed on the portable terminal while detailed information is printed by the printing apparatus. However, combining the portable terminal and printing apparatus faces several technical problems which must be solved.\nFor example, since the printing apparatus is large in size and is difficult to carry with the portable terminal, in order to print the detailed information, ordinarily, the printing apparatus that will be used is provided at a home or at an office. However, use of a fixed, particular printing apparatus would hardly allow information to be readily obtained anywhere. In order to achieve this, it is required that either the size of the printing apparatus be reduced to such an extent as to allow it to be carried, or an arbitrary printing apparatus be allowed to be used. The former is not feasible because it is technically very difficult to achieve with technology that is presently available. The latter is not feasible either because each printing apparatus requires a unique driver, thus requiring that the portable terminal be embedded with drivers for each of the printing apparatuses to be used.\nAccordingly, the present invention solves the above problems with the conventional art, and an object thereof is to provide a data output control apparatus which is suitable for allowing detailed information on a network to be readily obtained.\nIn order to achieve the above object, data output control apparatuses according to the present invention are provided, wherein detailed information on the network can be readily obtained compared with the conventional art. In particular, even if an output terminal is newly provided, setting of the data conversion device needs to be changed by the service provider only in relation to the new output terminal, allowing the new output terminal to be used without any change in setting by the user. Accordingly, the service provider will receive the advantages that setting works for newly provided output terminals are facilitated and that an information service with an adequate satisfaction can be provided to the user.\nIn accordance with the data output control apparatuses according to another aspect of the present invention, data conversion is executed at positions where the transmission load of the network or the processing load of the data conversion device is low; thus, irrespective of the transmission load of the network or the processing load of the data conversion device, the time required for receiving the service of output data is substantially constant. Accordingly, the service provider also receives the advantage that a print service with an even more comfortable printing environment can be provided to the user.\nIn accordance with the data output control apparatus according to another aspect of the present invention, detailed information on the network can be readily obtained compared with the conventional art. In particular, an optimal image is displayed on the portable terminal in accordance with the display functionality, allowing the user to grasp rather accurately an image of the data associated with the data print request as printed by the output terminal. In addition, by notifying the user of the preview data prior to the service of output data, the user is prevented from printing wrong data. Accordingly, the service provider receives the advantage that an information service with an adequate satisfaction can be provided to the user.\nIn accordance with the data output control apparatuses according to another aspect of the present invention, detailed information on the network can be readily obtained compared with the conventional art. In particular, even if an output terminal is newly provided, setting of the data conversion device needs to be changed by the service provider only in relation to the new output terminal, allowing the new output terminal to be used without any change in setting by the user. Accordingly, the service provider will receive the advantages that setting works for newly provided output terminals are facilitated and that an information service with an adequate satisfaction can be provided to the user.\nIn accordance with the data output control apparatuses according to another aspect of the present invention, data conversion is executed at positions where the transmission load of the network or the processing load of the data conversion device is low; thus, irrespective of the transmission load of the network or the processing load of the data conversion device, the time required to receive the service of output data is substantially constant. Accordingly, the service provider also receives the advantage that an output service with an even more comfortable output environment can be provided to the user."}
{"text": "The present invention relates to a drive nut a drive screw for a seat adjuster and in particular to a nut and screw of a vehicle seat assembly that provides ease of assembly and low noise during operation.\nPower seat adjuster are a popular option on many automotive vehicles and are frequently offered as standard equipment on higher priced vehicles. Such power seat adjusters are primarily used on the driver's seat and sometimes, on the front passenger seat of the vehicle to provide selective horizontal fore and aft and vertical movement of the seat as well to provide a comfortable seating position to suit each occupant's preference.\nSuch power seat adjusters typically carry a support frame, which supports the seat bottom and in some applications the seat back of the vehicle seat. The seat support frame is mounted on first and second, spaced track assemblies, each formed of an upper track, which is connected to the seat support frame and is slidably mounted on a lower track anchored to the vehicle floor. A drive mechanism typically includes a bi-directional electric motor, which rotates a pair of drive shafts extending outward from the motor to a gear assembly or box mounted on each upper track. In one arrangement, each gearbox rotates a lead screw extending longitudinally below each upper track. A drive block is fixedly mounted to each lower track and threadingly receives the lead screw to cause reciprocal, horizontal movement of the upper track and the attached seat support frame upon selective energization of the drive motor and the drive shafts.\nSimilar drive mechanisms are provided for vertical (up and down) adjustment of the seat support frame and, in a seat adjuster having a seat recliner adjustment mechanism, for angularly adjusting the position of the seat back with respect to the seat bottom. The vertical drive mechanism may also include separate front and rear edge seat adjustment drive mechanisms for selectively tilting the front rear edges of the seat bottom independent of each other as well as simultaneously to raise and lower the entire vehicle seat or, in some applications, only the seat bottom or cushion.\nEach vertical and recliner drive mechanism also includes a drive motor having a rotatable output shaft connected to a gear assembly either directly in the case of the front and rear vertical drive mechanisms or by means of two shafts extending a from single motor to separate gear assemblies mounted on each upper track or on the seat support frame in the case of a seat recliner drive mechanism.\nAs the various drive motors are mounted between the track assemblies, the drive motor output shafts and drive shafts connected thereto are typically oriented perpendicular to the axis of the associated lead screw. In order to transmit rotation and drive force between the output shaft of the drive motor and the drive shaft connected thereto to the 90 degree offset lead screw, each gear assembly box typically includes a worm gear, which is insert molded on one end of the associated lead screw. A worm is rotatably mounted in the gear box housing in meshing engagement with the worm gear and is connected to either one of the drive motor output shafts or to one of the drive shafts so as to rotate the worm gear and thereby the lead screw upon bi-directional energization of the selective drive motor. Thus, each drive motor requires a separate gear assembly for each lead assembly for each lead screw driven by the drive motor.\nThe rigid connection between the gears in the gear assembly, the drive shaft motor output shaft, the drive shaft, the lead screw and drive block also leads to additional problems during assembly and operation of a power seat adjuster. It is inevitable in the manufacture assembly of a mechanical mechanism, such as power seat adjuster, that dimensional conditions, such as concentricity, TIR, and linear discrepancies from nominal design dimensions, can and typically do occur. These dimensional conditions, without correction, can cause various problems in the operation of the power seat adjuster, such as poor breakaway from a stop position, slow operation of the power seat adjuster in extreme temperatures, excessive wear of the components of the power seat adjuster, uneven operation of the power seat adjuster.\nAnother aspect of a power seat adjuster, which is critical in the use of a power seat adjuster, particularly during excessive forces generated during a vehicle collision, is the requirement for seat integrity. As the occupant's seat belt is frequently attached by a seat belt buckle mounting bracket directly to one of the upper tracks of a power seat adjuster to enable the seat belt to move fore and aft with the upper track, any forces exerted on the occupant during a vehicle collision are transmitted directly though the seat belt buckle mounting bracket to the upper track. These forces cause the upper track to move upward with respect to the lower track and, in the event of excessive force, could lead to a complete pullout or separation of the upper track and the attached seat from the lower track. In order to prevent track separation, seat integrity requirements have necessitated the design of the components of a power seat adjuster to resist track separation.\nSeveral attempts have been made to provide adapters that meet the above requirements. For example, in U.S. Pat. Nos. 5,172,601, 5,467,957, and 5,575,531, a drive nut with a pair of opposing flat sided legs with rounded edges that engages a bracket having wings with a pair of opposing open ended keyhole slots is disclosed. The slots have an arcuate portion in communication with a pair of opposing flat sides. The legs are inserted into slots and then rotated so that the rounded edges of the legs engage the arcuate portion of the slots. However, the load carrying capability of the wings is reduced because of the size of the opening in the slots required to insert the legs. As a result, these designs are not always able to resist separation forces on the seat and are prone to structural failure by allowing the slots to bend and the legs to separate from the slots causing separation of the seat.\nAnother attempt to solve the above problem is shown in U.S. Pat. No. 5,860,319. This design uses a round block drive device to fit into a housing. A resilient member is used to cushion or isolate the drive block from the housing. The resilient member co-acts with holes formed in the side walls of the housing to permit relative motion of the threaded interconnected lead screw and drive block with the housing. This design is expensive and difficult to assembly and has not been widely used in vehicle seat applications.\nThus, the vehicle seat assemblies and, in particular seat assemblies having power adjustment capability, it is desirable to provide an adjuster, which produces low noise, emissions during operation prevents seat separation and is simple and inexpensive to make. Furthermore, it is desirable to provide an adjuster that can be easily assembled to reduce assembly time and errors.\nThe present invention seeks to solve most of the above problems with a simple, quiet, easy to assemble drive nut that is structurally superior to known present designs.\nA vehicle seat linkage assembly including a threaded rod defining a longitudinal axis. The linkage includes a first link having a portion forming an enclosed aperture. A second link is adjacent but spaced away from the first link. The second link having a portion forming a slot. The slot has a circular portion. The circular portion and the enclosed aperture define a lateral axis extending between the first link and the second link. Additionally, a drive nut threadably engages the rod. The drive nut has one lateral projection, another lateral projection opposite the one projection and a portion engaging the rod. The drive nut and rod are moveable so that the longitudinal axis is positioned relative to the lateral axis to insert the one projection into the slot and to position the another projection adjacent but spaced away from the aperture while the one lateral projection remains in the slot. The another projection has a width substantially the same as the base of the aperture. Thereafter moving the another projection relative to insert the another projection into the aperture while the one lateral projection remains in the slot.\nThe present invention is advantageous in that the drive nut engages a slot on one side and an aperture on the other side so that the drive nut is prevented from being withdrawn from the slot and the aperture. Furthermore, the drive nut is easy to assemble, inexpensive to make and produces law noise.\nFrom the foregoing disclosure and the following more detailed description of the various preferred embodiments, it will be apparent to those skilled in the art that the present invention provides a significant advance in the technology and art of vehicle seat linkage assemblies. Additional features and advantages of various preferred embodiments will be better understood in view of the detailed description provided below.\nIt should be understood that the appended drawings are not necessarily to scale, presenting a somewhat simplified representation of various preferred features illustrative of the basic principles of the invention. The specific design features of the seat adjuster assembly as disclosed here, including, for example, specific dimensions of the drive nut and its associated bracket, will be determined in part by the particular intended application and use environment. Certain feature of the illustrated embodiments have been enlarged or distorted relative to others to facilitate visualization and clear understanding. In particular, thin features may be thickened, for example, for clarity of illustration. All references to direction and position, unless otherwise indicated, refer to the orientation of the drive member in the seat structure illustrated in the drawings."}
{"text": "As a result of a continuous increase in the integration degree and reliability of an integrated circuit, an analog circuit and a digital circuit are integrated and developed in one chip. Also, an analog circuit and a system play an important role in realization and application of a high-density Integrated Circuit (IC) technology. For example, nearly all high-density IC systems employ an amplifier, a filter, a detector, a comparator, etc.\nFIG. 1 illustrates an analog filter structure where analog variable amplifiers are configured in two stages according to the related art.\nReferring to FIG. 1, an analog variable amplifier 100 of a first stage includes an operational amplifier and two variable resistors R1 and R2. Depending on the implementation, the analog variable amplifier 100 of the first stage may include one variable resistor and one variable capacitor instead of the operational amplifier and the two variable resistors R1 and R2.\nLikewise, an analog variable amplifier 102 of a second stage includes an operational amplifier and two variable resistors R3 and R4. Depending on the implementation, the analog variable amplifier 102 of the second stage may include one variable resistor and one variable capacitor instead of the operational amplifier and the two variable resistors R3 and R4.\nIn each analog variable amplifier stage, a gain is determined by a ratio of the input resistors R1 and R3 to the feedback resistors R2 and R4. Assuming a transfer function of the operational amplifier is A(s), the gain is expressed below by Equation (1). In an ideal operational amplifier, a gain has a value of R2/R1 or R4/R3 in its infinity. Also, the gain of the entire circuit is represented as product (G1×G2) of gains at respective stages. Here, G1 is the gain of the analog variable amplifier 100 of the first stage, and G2 is the gain of the analog variable amplifier 102 of the second stage.\n                                          G            1                    =                                                    R                2                            ⁢                              A                ⁡                                  (                  s                  )                                                                                    R                1                            +                              R                2                            -                                                R                  1                                ⁢                                  A                  ⁡                                      (                    s                    )                                                                                      ⁢                                  ⁢                              G            2                    =                                                    R                4                            ⁢                              A                ⁡                                  (                  s                  )                                                                                    R                3                            +                              R                4                            -                                                R                  3                                ⁢                                  A                  ⁡                                      (                    s                    )                                                                                                          Equation        ⁢                                  ⁢                  (          1          )                    \nHere, the variable gain amplifier is configured by controlling the values of the input resistors R1 and R3 or the feedback resistors R2 and R4. In the analog variable amplifier 100 of the first stage, when R1 is greater than R2, the variable gain amplifier performs attenuation. When R1 is less than R2, the variable gain amplifier performs amplification. Likewise, in the analog variable amplifier 102 of the second stage, when R3 is greater than R4, the variable gain amplifier performs attenuation. When R3 is less than R4, the variable gain amplifier performs amplification.\nMeanwhile, in a case of configuring an analog amplifier requiring a range from −40 dB or less to 0 dB or more, an analog filter processes a wide operation range via two or more analog amplifier stages. For example, a first analog amplifier stage processes −20 dB or less ˜0 dB or more, and a second analog amplifier stage processes −20 dB or less ˜0 dB or more, so that the entire gain becomes −40 dB or less ˜0 dB or more.\nHowever, an analog filter including a plurality of amplifier stages consumes twice more power compared to the case where one amplifier stage amplifies an analog signal, and requires twice more area, so that manufacturing costs increase.\nTherefore, there is a need for an analog amplifier including a single amplifier stage.\nThe above information is presented as background information only to assist with an understanding of the present disclosure. No determination has been made, and no assertion is made, as to whether any of the above might be applicable as prior art with regard to the present disclosure."}
{"text": "In recent years, an electronic device (compound semiconductor device), in which a GaN layer and an AlGaN layer are disposed sequentially above a substrate and the GaN layer is used as an electron transit layer, has been developed actively. As for an example of such compound semiconductor devices, a GaN based high electron mobility transistor (HEMT) is mentioned. In the GaN based HEMT, a high-concentration two-dimensional electron gas (2DEG) generated at a heterojunction interface between AlGaN and GaN is utilized.\nThe band gap of GaN is 3.4 eV and is larger than the band gap of Si (1.1 eV) and the band gap of GaAs (1.4 eV). That is, GaN has a high breakdown field strength. Furthermore, GaN has also a large saturation electron velocity. Therefore, GaN is a very promising material for a compound semiconductor device capable of high-voltage operation and production of high output, for example, a material for a semiconductor device for a power supply. Consequently, the compound semiconductor device by using the GaN based compound semiconductor device is expected as a high-breakdown voltage power device for a high-efficiency switching element, an electric car, and the like.\nRegarding the GaN based HEMT, the material for a gate electrode is different from the material for the source electrode and the drain electrode. Therefore, the gate electrode is formed by a process different from that of the source electrode and the drain electrode. The gate electrode, the source electrode, and the drain electrode are formed by, for example, a lift-off method. That is, in formation of the electrode, formation of a resist pattern, formation of an electrode material, and removal of the resist pattern are performed. Meanwhile, in production of a GaN based HEMT, recesses or opening portions may be formed in regions to be provided with the gate electrode, the source electrode, and the drain electrode of the compound semiconductor layer. In this case, a compound semiconductor layer is etched by using a resist pattern, so as to form a recess or an opening portion, and thereafter, the resist is removed. Furthermore, after the source electrode and the drain electrode are formed, the source electrode and the drain electrode may be covered with a passivation film, and before formation of the gate electrode, the passivation film may be dry-etched by using a resist pattern. In this case, the resist is removed after the dry etching.\nIn these methods, although the resist is removed, residues of the resist may remain. Moreover, after the dry etching of the passivation film, etching residues may remain. In this regard, these residues may not be removed sufficiently even by cleaning with an organic material. This is because of effects of alteration and the like during post baking after development performed in formation of the resist pattern and dry etching.\nMeanwhile, when the residues are removed by an acid treatment through the use of, for example, a mixture of sulfuric acid and aqueous hydrogen peroxide, an electrode exposed at that time is damaged. For example, in the case where a recess for the gate electrode is formed after the source electrode and the drain electrode are formed, the source electrode and the drain electrode are exposed during formation of the recess. Therefore, if the above-described mixture is used, the source electrode and the drain electrode are damaged. Alternatively, even when the electrodes are not damaged, the exposed surface of the compound semiconductor layer may undergo oxidation, surface roughening, and the like, so that the characteristics may be changed.\nThen, in the case where the above-described residues are present, an increase in leakage current, fluctuations in threshold voltage due to trap of charge, and the like occur, so that the yield is reduced significantly.\nThe followings are reference documents:\n[Document 1] Japanese Laid-open Patent publication No. 2003-167360 and\n[Document 2] Japanese Laid-open Patent publication No. 2007-128038."}
{"text": "Intracardiac and intravascular procedures, whether performed percutaneously or in an open, surgical, fashion, may liberate particulate debris. Such debris, once free in the vascular system, may cause complications including vascular occlusion, end-organ ischemia, stroke, and heart attack. Ideally, this debris is filtered from the vascular system before it can travel to distal organ beds.\nUsing known filter mechanisms deployed in the arterial system, debris is captured during systole. There is a danger, however, that such debris may escape the filter mechanism during diastole or during filter removal. Apparatus and methods to reduce debris escape during diastole or during filter removal may be desirable to reduce embolic complications."}
{"text": "The present invention relates to the field of computing. More particularly, it relates to a method for performing analysis of data.\nOn-Line Analytical Processing (OLAP) is a technology that allows end-users to retrieve data efficiently from one or more data warehouses for decision-support purposes. A data warehouse can be defined as an on-line repository of enterprise data used to support decision making OLAP transforms data warehouse data into strategic information. Data warehouses tend to be extremely largexe2x80x94hundreds of gigabytes to terabytes in size. Information in a data warehouse is usually multidimensional in nature, and end-users require the capability to view that data from a variety of perspectives.\nIn this milieu, aggregated and summarised data are much more important than detailed records. Queries tend to be complex and ad-hoc, often requiring computationally expensive operations such as joins and aggregation. Further complicating this situation is the fact that such queries must be performed on tables that potentially have millions of records. Moreover, the results have to be delivered interactively to the end-user who expects rapid response times.\nOLAP enables end-users to gain insight into data through fast, interactive views of information. As well as the ability to answer xe2x80x9cwho?xe2x80x9d and xe2x80x9cwhat?xe2x80x9d questions they have ability to answer xe2x80x9cwhat if?xe2x80x9d and xe2x80x9cwhy?xe2x80x9d\nGiven these characteristics and the consequent processing requirements, it is clear that there must be ongoing development emphasis in OLAP systems in the area of query processing and response times. Processing is critical in such an environment, and any measures that improve the efficiency of processing are welcome, as they also serve to reduce the response times for the end-user.\nAccording to the invention, there is provided a computer-based method for analyzing data comprising the steps of: at a central application platform: a) receiving from an end-user a query; b) dividing said query into a number of partial queries; c) designating said partial queries as simple or complex as appropriate; c) sending said simple queries to one or more OLAP data sources for evaluation; d) emulating the evaluation of partial queries not of the simple type; e) receiving from said OLAP data sources the results of the evaluation from step c); and f) combining the results of said evaluations and said emulation; and at one or more OLAP data sources; g) receiving from a central application platform one or more requests to evaluate simple queries; and h) transferring to said central application platform the results of evaluating said simple queries."}
{"text": "1. Field of the Invention\nThe invention relates in general to a signal input apparatus, and more particularly to a signal input apparatus having alternative power supply while the battery is exhausted\n2. Description of the Related Art\nComputer mouse is a signal input apparatus of a computer, especially being a pointing device for controlling the displacement of the cursor on screen. Conventionally, computer mouse is classified as an optical mouse and a mechanical mouse. The mechanical mouse controls the displacement of the cursor according to the input signal produced in accordance with the relative position of the rolling ball. Compared to the mechanical mouse, the optical mouse detects the relative movement of the mouse by an optical sensor and the optical sensor includes a light emitting diode (LED), lens and sensing IC. The sensing IC is preferably an IC of Agilent Technologies Ltd, the model No. of which is ADNS-2610 (referring to the U.S. Pat. No. 6,433,780).\nMoreover, it brings more convenience to users that the computer mouse communicates with computer by wireless technique. When a computer and a projector are common applications in business, a wireless computer mouse allows a speaker to control the computer and move simultaneously during presentation.\nHowever, a wireless computer mouse is driven by power of a battery, and the wireless computer mouse cannot be driven when the power of the battery is exhausted. Also, a pointing stick or a laser pen is required for the speaker to guide listeners' sight during presentation, but it is not easy for the speaker to hold a wireless computer mouse and a laser pen in one hand."}
{"text": "1. Field of the Invention:\nThe present invention relates to a chain which is wound around a V-pulley and is suitably used for a continuously variable transmission of an automobile or other transmissions in general.\n2. Related Art:\nIn recent years, a continuously variable transmission system comprising a pair of V-pulleys and an endless belt or chain has come to be used for a continuously variable transmission which is coupled with an engine of such as an automobile having a relatively large torque. Specifically, this continuously variable transmission system is arranged as follows. A pair of V-pulleys in which an interval between their inner walls can be varied are mounted on a drive shaft and a driven shaft, respectively, an endless belt or chain is trained between the two V-pulleys, and the interval between the inner walls of each V-pulley is varied by, for instance, hydraulic means in response to speed change commands. This causes the winding radius of the endless belt or chain trained between the V-pulleys to be varied relatively and continuously, thereby effecting a continuously variable transmission between the drive shaft and the driven shaft.\nAn endless belt described above is disclosed in Japanese Pat. Publication No. 6783/1980. This is a metal V-belt which is arranged such that metal blocks with a thickness of about 2 mm each having a surface of contact with an inner wall of a V-pulley are fitted with practically no gap with a hoop formed by superposing a plurality of thin metal sheets with no gap therebetween. The transmission of a driving force in this V-belt is effected as follows: The torque of the drive shaft-side V-pulley is transmitted to the blocks via a frictional force, and the driving force between the blocks is transmitted to the preceding blocks through a pushing force. Meanwhile, the driving force at the driven shaft-side V-pulley is converted again to torque through a frictional force. That is, the driving force does not substantially act on the hoop, and the hoop serves as a binding member for preventing the blocks from becoming separated. Accordingly, with this V-belt, no fracture of the hoop takes place, and it has been considered that it excels in terms of durability and noise. However, since the driving force is transmitted by the pushing force exerted among the blocks, the V-belt on the driven shaft-side V-pulley acts to escape in the direction of the outer periphery of the V-pulley. Consequently, there has been a drawback in that the friction between the V-belt and the V-pulley is made disadvantageously small, thereby causing a decline in the driving force-transmitting capabilities. In addition, it has been very difficult to fabricate a hoop formed by superposing a plurality of thin metal sheets without leaving any gap therebetween, and much manpower and time have been required in the fabrication thereof.\nAccordingly, a chain for a V-pulley has been proposed in which blocks coming into contact with the V-pulley are fixedly secured to an ordinary chain. According to this chain for a V-pulley, since the driving force is transmitted through a pulling force of the chain, it is possible to overcome the problem of a decline in the driving force transmitting capabilities, which is the drawback of the above-described V-belt. Such chains for V-pulleys are disclosed in, for instance, Japanese Pat. Laid-Open Nos. 226729/1984, 226730/1984, 26831/1985, and 143243/1985. However, these prior arts are generally concerned with techniques for attaching blocks to the chain, or forming the blocks on the link plates or on the opposite ends of the connecting pins. Thus, the end faces of the blocks which are brought into contact with the inner wall of V-pulley have generally been formed into flat surfaces with no particular concern being paid to the configuration of the end faces of the blocks.\nAs an example in which the configuration of the end face of the block is made not flat, it is possible to cite a technique disclosed in, for instance, Japanese Utility Model Laid-Open No. 196850/1988. In this technique, a cylindrical body having an axis in the longitudinal direction of a chain abuts at a peripheral surface thereof against an inner wall of a V-pulley. However, the inner wall of the V-pulley is conically shaped, and, with this technique, the blocks and the V-pulley can theoretically be only brought into point contact with each other, so that it is considered that it is disadvantageous in terms of the driving force-transmitting capabilities.\nIn this respect, if the end face of the block is formed with a flat surface having the same inclination as that of the inner wall of the V-pulley (i.e., half the apex angle), it is possible for the end face of the block and the inner wall of the V-pulley to be theoretically brought into contact with each other at a linear line of contact, and this arrangement is advantageous in the transmission of the driving force- However, when the block starts to be engaged with the V-pulley, the block is brought into contact with the V-pulley at a front edge thereof, and when it begins to be disengaged with the V-pulley, the block is brought into contact with the V-pulley at a rear edge thereof. In other words, at the beginning of engagement and disengagement, the block and the V-pulley are not brought into smooth linear contact with each other and the block comes into contact with the V-pulley at its edge, thereby causing a decline in the driving force-transmitting capabilities. In addition, during a period from the time when the block enters the V-pulley and is engaged therewith and until the block is disengaged therefrom, the block rotates back and forth about an axis parallel to the axis of the connecting pin. Since this rotation of the block is effected under high pressure, rotational friction occurs between the block and the V-pulley, causing abrasion of both the block and the V-pulley and a loss in the driving force- In addition, since the portion of the end face of the block which becomes worn is fixed, it is impossible to obtain a chain for a V-pulley having an extended life. Furthermore, since the chain disposed between the two V-pulleys vibrates in the axial direction of the connecting pin, the block collides against the inner wall of the V-pulley at the time of entering the V-pulley, thereby generating noises due to this collision.\nTo solve these problems, the present inventors proposed a technique (Japanese Pat. Laid-Open No. 219937/1988) in which the block is connected to the tip of the connecting pin in the form of a universal joint without the block being fixed to the link plate, and the end face of the block is made rotatable and swingable to a certain degree about the axis of the connecting pin while maintaining the same inclination as that of the inner wall of the V-pulley. According to this technique, the block is rotatable, the rotational friction caused by the rotation of the block occurring at the time of entrance into the V-pulley does not occur, and the block is swingable. Therefore, it had been expected that the colliding sound occurring when the block collides against the inner wall of the V-pulley during entrance into the V-pulley would be reduced. However, it was not possible to attain a reduction of the noises to the extent that had been expected. As the reason for this problem, it can be conjectured that since the end face of the block was formed into a flat surface, the block and the V-pulley could not maintain smooth linear contact."}
{"text": "Generally described, computing devices utilize a communication network, or a series of communication networks, to exchange data. Companies and organizations operate computer networks that interconnect a number of computing devices to support operations or provide services to third parties. The computing systems can be located in a single geographic location or located in multiple, distinct geographic locations (e.g., interconnected via private or public communication networks). Specifically, data centers or data processing centers, herein generally referred to as a “data center,” may include a number of interconnected computing systems to provide computing resources to users of the data center. The data centers may be private data centers operated on behalf of an organization or public data centers operated on behalf, or for the benefit of, the general public.\nTo facilitate increased utilization of data center resources, virtualization technologies may allow a single physical computing device to host one or more instances of virtual machines that appear and operate as independent computing devices to users of a data center. Each single physical computing device can be generally referred to as a host computing device. With virtualization, the single physical computing device can create, maintain, delete, or otherwise manage virtual machines in a dynamic matter. In turn, users can request computer resources from a data center, including single computing devices or a configuration of networked computing devices, and be provided with varying numbers of virtual machine resources.\nIn conjunction with the utilization of virtualization technologies, data centers can physically organize sets of host computing devices to allow the host computing devices to share computing device resources, such as power or communication network connectivity. Such physical organization can correspond to physical racks in which the hosting computing devices are mounted, generally referred to as racks of host computing devices. As the number of racks of host computing devices increases, service providers associated with data centers have difficulty distinguishing between errors or faults associated with individual host computing devices, shared resources associated with a particular rack or distributed components utilized to manage the host computing devices. Additionally, in some scenarios, service providers may maintain multiple data centers. In such scenarios, the service providers would benefit from the determination of the availability of multiple data centers and the inter-data center synchronization of state information."}
{"text": "The present disclosure relates generally to the field of a system and method for receiving an audible input in a vehicle. More specifically, the disclosure relates to a steering wheel having at least one microphone and configured to process audible input.\nA communication steering wheel device has been developed in which a variety of informational messages can be displayed to the driver primarily through icons. These messages include a variety of warning messages, informational messages about upcoming landmarks, navigation commands, etc. For example, see US Published Application US20060070795 and WIPO Publications W02006076903 and W02006076904 (all three applications incorporated by reference herein).\nA vehicle may include microphones to provide or control additional features (e.g., hands free telephone use, navigation control, etc.) using voice recognition algorithms. Much attention has been placed in the industry to improve and refine voice recognition but significant challenges remain to make such systems sufficiently reliable for automotive use. For example, detection of the driver voice can be obscured by other voices and sources of noise in and outside of the car. Currently, an input device (microphone) may be included within the instrument panel (IP) or headliner of the vehicle. There is a need for a system and method for filtering out road noise, car noise, other occupants, and other noise for voice recognition."}
{"text": "1. Field of Invention\nThe present invention relates to a radio frequency identification (RFID) reader using a plurality of antennas for reading data from Tags, and more particularly to a RFID reader using a plurality of antennas applied to multi-dimensional objects, such as three-dimensional jigsaw with multimedia sound and light effect, to distinguish the orientation of the three-dimensional jigsaw unit between right and wrong.\n2. Description of Related Arts\nRadio Frequency Identification (RFID) systems are used for identification and/or tracking of equipment, inventory, or living things. RFID systems are radio communication systems that communicate between a radio transceiver, called a RFID reader, and a number of inexpensive devices called Tags or transponders. In RFID systems, the RFID reader communicates to the Tags using modulated radio signals, and the Tags respond with modulated radio signals. Conventional RFID systems are designed a) to identify an object passing into range of the RFID reader, and b) to store data onto the Tag and then retrieve that data from the Tag at a later time in order to manage inventory or perform some other useful application. In this application, the transmit and receive antennas have the same bandwidth. In fact, transmit and receive frequently share the same antenna, using a circulator to separate the transmit and receive paths.\nAs described in the U.S. patent application Ser. No. 11/062,815, filed on Feb. 23, 2005, a radio frequency identification (RFID) reader comprising a plurality of antennas could be applied to multi-dimensional objects, such as three-dimensional jigsaw with multimedia sound and light effect, Monopoly board game, . . . etc. However, in this prior art each of antennas in the RFID reader is positioned in a central region of each plane of multi-dimensional objects such as three-dimensional jigsaw. In order for the RFID reader to sense the identifier, the polarizations of the electromagnetic fields generated by the antennas in the RFID reader and Tag must be aligned. In most real-world implementations, however, the orientation of the Tag, and, hence, the polarization of the field generated by the Tag's antenna, is unknown. Therefore, the RFID reader could not distinguish the orientation of the three-dimensional jigsaw unit between right and wrong because the modulated radio signals are non-direction signals. In other words, the RFID reader could not distinguish the orientation of the three-dimensional jigsaw unit between right and wrong even though the three-dimensional jigsaw unit aims at the right position. Therefore, there is a need in the art for a RFID reader that can read Tags and also distinguish the Tags' orientations."}
{"text": "Threaded unions are used to provide fluid-tight joints in fluid conduits. Threaded unions are held together by a threaded nut that is tightened to a required torque using a wrench or a hammer. In the oil industry, threaded unions are generally constructed using “wing nuts” and are commonly called “hammer unions” or “hammer lug unions”. Hammer unions are designed and manufactured in accordance with the specifications stipulated by the American Petroleum Institute in API 6A entitled “Specification for Wellhead and Christmas Tree Equipment”. Hammer unions are usually available in a variety of sizes (1″ to 10″) and a variety of pressure ratings (1000 psi to over 20,000 psi).\nAs illustrated in FIG. 1, a typical prior-art hammer union 10 includes a first subcomponent 12 and a second subcomponent 14. The first subcomponent has an inner lateral surface 13 which abuts an outer lateral surface 15 of the second subcomponent. The first subcomponent has at least one annular groove 16 for receiving an elastomeric O-ring 18. The annular groove is located at an interface of the inner and outer lateral surfaces. The second subcomponent is secured to the first subcomponent by a wing nut 20.\nThe wing nut 20 has box threads 22 that engage pin threads of the first subcomponent. The wing nut 20 further includes a plurality of lugs 24 that extend radially from a main body of the wing nut 20. The lugs have impact surfaces 25 which may be impact-torqued using a hammer or mallet (not shown) to tighten or loosen the wing nut 20.\nThe wing nut 20 also has an annular top wall 26 which abuts a radial flange 28 of the second subcomponent. When torque is applied to the wing nut, the annular top wall 26 is forced downwardly on the radial flange 28, thereby locking together the second subcomponent and the first subcomponent.\nThe hammer union 10 is shown with an elastomeric O-ring 18 having a circular cross-section. In the prior art, the O-rings are made of materials such as rubber, nylon, polyurethane and Teflon™.\nDepending on the nature and shape of the first and second subcomponents, a different type of elastomeric seal may be used. For example, lip seals and flat gaskets may be used instead of the O-ring shown in FIG. 1. For sour service wells, the lips seals and O-rings are typically made of nitrile rubber (NBR) and fluroelastomers (FPM). Elastomeric seals may also be energized using stainless steel garter springs to prevent the seal from being damaged by extrusion.\nFIG. 2 is a schematic cross-sectional view of a prior art hammer union manufactured by FMC/Weco that utilizes a lip seal. The hammer union shown in FIG. 2 is very much like the hammer union shown in FIG. 1, with the exception that the high-pressure seal is provided by the lip seal 17, which includes a peripheral lip 19 that is received in a groove 21 at a bottom of a seal seat 23 in the first subcomponent 12. The lip seal 17 is made of a rubber composition, and the hammer union is rated for up to 15,000 psi of fluid pressure.\nOne substantial disadvantage of prior-art hammer unions is that their elastomeric seals are vulnerable to the extreme temperatures generated by fire. In the event that a fire erupts in the well or at the wellhead, the elastomeric seal in the hammer union may leak or fail completely. This permits hydrocarbons to escape to the atmosphere, which may exacerbate the fire.\nTo the best of the applicant's knowledge, a prior art hammer union with a metal-to-metal seal for providing a high-pressure, fluid-tight seal has never been designed or manufactured. While metal seals and metallic gaskets are known in the art (e.g. U.S. Pat. No. 4,832,381 (Boulton) entitled “Seal”, U.S. Pat. No. 5,257,792 (Putch et al.) entitled “Well Head Metal Seal”, U.S. Pat. No. 4,056,272 (Morrill) entitled “Seal”, and U.S. Pat. No. 1,825,962 (Laird) entitled “Gasket”) but each of the above describes a seal for a flanged union. While flanged unions are widely used in well trees, they are relatively expensive to construct and time consuming to assemble in the field.\nIt is well known in the art that there is increasing pressure on the oil industry to produce hydrocarbons at a lower cost. Consequently, an interest has developed in utilizing wellhead equipment that is less expensive to construct and is more quickly assembled than prior art flanged well tree components. Threaded unions provide a good alternative to flanged unions from a cost standpoint because they are faster to assemble and less expensive to construct. However, due to safety concerns related to the lack of a metal-to-metal seal, use of threaded unions for well tree components has not been endorsed.\nTherefore, there exists a need for a threaded union with a metal-to-metal seal."}
{"text": "Tablets comprising a cellulose ether are known to the pharmaceutical drug delivery art. For example, tablets containing the cellulose ether hydroxypropylmethycellulose are known in U.S. Pat. Nos. 3,870,790; 4,140,755; 4,167,588; 4,226,849; 4,259,314; 4,357,469; 4,369,172; 4,389,393 and 4,540,566.\nThe tablets known to the prior art using the hydroxypropyl-methylcellulose ether often have certain disadvantages associated with their structure and with their use. For example, the mechanical integrity of some prior art tablets frequently is insufficient to provide both a sustained and a rate controlled release of a drug over a prolonged period of time in a moving fluid environment of use. The prior art tablets often exhibit insufficient mechanical integrity, that is the cohesive ability to stay together in a moving fluid environment such as the gastrointestinal tract, without prematurely breaking-up and without prematurely releasing all of its drug content. The above-mentioned desirable properties are not readily apparent in the prior art tablets, which appear to undergo substantial disintegration in a short time span, usually less than eight hours in a fluid environment of use.\nAnother disadvantage associated with the prior art tablets is that they exhibit an unwanted, variable, and difficult to reproduce a rate of release pattern. For example, prior art tablets comprising a small amount of a cellulose ether exhibit this behavior, such as by tablets consisting of less than five weight percent of a hydroxypropylmethylcellulose having a number average molecular weight greater than 50,000. The presence of the small amount of this high molecular weight polymeric ether in the tablet masks the release characteristic of other polymeric ethers in the tablets resulting in an erratic release pattern which is difficult to reproduce from batch to batch and from tablet to tablet.\nStill other unacceptable disadvantages associated with the prior art tablets are that the tablets during their shelf-life can exhibit an unpredictable change in their release-rate characteristics; the prior art tablet when tested in an in vitro test that substantially reproduces the in vivo environment of the gastrointestinal tract often releases the drug at a greater rate of release in vivo than in vitro, which difference can be attributed to a premature disintegration of the prior art tablet; and the prior art tablet in a high fluid shear environment releases its drug too quickly, usually in less than six hours and these tablets therefore are not adapted to prolonged release.\nThus, in the light of the above presentation it will be appreciated by those versed in the dispensing art, that if a novel dosage form is made available to the medical and the pharmaceutical arts for dispensing a difficult to deliver drug free of the tribulation known to the prior art, such a dosage form would have a definite use and would also be a valuable contribution to the dispensing art. It will be further appreciated by those versed in the dispensing art that if a dosage form can be provided that (a) possesses a desirable rate of release and mechanical properties for dispensing a drug over a prolonged period of time, and which dosage form (b) can be manufactured at an economical cost, such a dosage form would have a positive and a practical value and it would also represent an advancement in the dispensing arts."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device for power supply, and, in particular, to a DC-DC converter which is a kind of a switching-type constant-voltage power supply in which, when an externally supplied voltage given by an external supply power source such as a battery or the like is converted into a predetermined constant voltage, and, then, is outputted, determination is made as to whether a condition of a load to which the constant voltage is provided is a heavy load condition or a light load condition, and switching of a switching transistor is controlled in accordance with the load condition so that an improvement in power-supply efficiency is achieved at the time of a light load.\n2. Description of the Prior Art\nAs such a kind of DC-DC converter in the prior art, there is one which is disclosed in Japanese Laid-Open Patent Application No. 7-322608 (see FIG. 1), for example.\nSpecifically, a DC-DC converter 2A in the prior art, shown in FIG. 1, supplies power, via a common power-supply line 18A, to a heavy-load circuit 12A such as a DC-motor driving circuit or the like and a light-load circuit 15A such as an IC circuit generating a control signal or the like. The DC-DC converter 2A has a switching circuit including a transistor 5A which is inserted in series or in parallel to the heavy-load circuit and light-load circuit and performing switching control of the transistor 5A for the power-supply line 18A, a control circuit 6A which controls periods of switching of the transistor 5A through pulse-width modulation (PWM) control or pulse-frequency modulation (PFM) control according to a voltage V001 at a load side so that this voltage V001 is fixed, and an oscillation circuit 7A which turns on and turns off the transistor 5A in a predetermined period.\nIn the DC-DC converter having the above-described circuit arrangement, a control unit 20A (MCU) causes a battery voltage Vcc of a battery 16A to be provided to the control circuit 6A via a switch 8A, causing the control circuit 6A to operate so that the control circuit 6A performs the pulse-width modulation control, when power is supplied to the heavy-load circuit 12A. The control unit 20A causes the battery voltage Vcc of the battery 16A not to be provided to the control circuit 6A, causing the control circuit 6A to stop operating, and, simultaneously, causing the battery voltage Vcc of the battery 16A to be provided to the oscillation circuit 7A via the switch 8A, causing the oscillation circuit 7A to operate so that the oscillation circuit 7A performs the pulse-frequency modulation control, when power is supplied only to the light load circuit 15A.\nHowever, in this DC-DC converter 2A in the prior art, because switching is performed such that the pulse-width modulation control or the pulse-frequency modulation control is used according to the load condition, although the efficiency is improved when the switching is performed such that the pulse-frequency modulation control is used at the time of light load, the load-side output direct-current voltage problematically includes a noise having a random frequency."}
{"text": "This invention relates to fuel cells and, in particular, to improvements in the catalyst used in internal reforming fuel cells.\nThe fuel cell is emerging as a viable device for efficiently converting chemical energy of a fuel directly to electrical energy. In a fuel cell, energy conversion is accomplished via electrochemical reactions involving a fuel and an oxidant taking place within the fuel cell. Although the fuel cell is an efficient converter of chemical energy to direct electrical power, the part of the fuel energy which is not converted to electricity in the fuel cell is available as heat. This by-product heat is removed to achieve isothermal operation. Water is another by-product of the fuel cell reactions.\nThe fuel used in practical and commercially viable fuel cells is hydrogen. Hydrogen is produced by steam reforming of a hydrocarbon fuel in an endothermic reforming reaction. Therefore, it requires a supply of both water and heat from an external source. In a conventional fuel cell power plant, an external reformer is used (see, e.g., U.S. Pat. No. 3,909,299) to produce hydrogen for the fuel cell. To improve overall efficiency of energy conversion, the fuel cell by-product heat is utilized in the reformer. This requires heat exchange equipment and hot piping.\nBecause the fuel cell reaction produces water and heat and the reforming reaction consumes water and heat, a technique for combining both of these reactions within the fuel cell, called internal reforming, has been proposed (see, e.g., U.S. Pat. Nos. 4,182,795 and 4,877,693). In the internal reforming fuel cell, the reforming reaction is carried out in situ in the fuel cell anode compartment so that the fuel cell produced water and heat is directly available at the reforming reaction site. This requires that the fuel cell anode compartment be loaded with an appropriate reforming catalyst.\nThe fuel cell anode compartment is usually a planar structure. This structure typically comprises a separator plate for isolating fuel from the neighboring cell oxidant stream, an anode for providing fuel cell reaction sites, and a current collector for conducting of electrons from the anode to the cathode of a next neighboring cell and for providing a flow path for the gaseous fuel stream.\nAs indicated above, the anode compartment of the internally reforming fuel cell is loaded with the reforming catalyst. Various methods for loading reforming catalyst in the fuel cell anode compartment have been used (see, e.g., U.S. Pat. No. 4,788,110). The reforming catalysts are usually available in different shapes such as tablet, pellet, rod, ring and spherical. In the conventional methods, the reforming catalyst of the pellet and/or the tablet shape are loaded into the pockets of the corrugated current collector.\nThe aforesaid catalyst loading methods suffer from several drawbacks. Uniform loading of the catalysts in the direction normal to the gas flow as well as in the direction of the gas flow is difficult to achieve resulting in fuel flow maldistribution, in the fuel cell. Handling of small catalyst particles during assembly is cumbersome and is difficult to automate and, therefore, is not cost effective. Also, cell-to-cell loading uniformity is difficult to achieve and catalyst shifting and spilling during assembly, handling, transportation, and during operation may be difficult to avoid.\nFurthermore, it is important that any contact between the catalyst and the anode be avoided during operation to prevent catalyst deactivation by electrolyte adsorption. In the most advanced designs to date, this has been achieved only if the fuel cells are situated vertically. Angular or the horizontal positions are not permitted because the catalyst particles in an inclined cell will shift and/or may come in contact with the anode and adsorb electrolyte resulting in deactivation and loss of reforming activity.\nIt is, therefore, an object of the present invention to provide a catalyst assembly for an internal reforming fuel cell which avoids these disadvantages.\nIt is a further object of the present invention to provide a catalyst assembly for an internal reforming fuel cell in which the catalyst assembly is easily assembled in the fuel cell.\nIt is yet a further object of the present invention to provide a catalyst assembly for an internal reforming fuel cell in which the catalyst assembly has minimum displacement and permits vertical and horizontal orientations of the fuel cell."}
{"text": "One type of semiconductor package is referred to as a BGA package. BGA packages were developed to provide a higher lead count, and a smaller foot print, than conventional plastic or ceramic semiconductor packages. A BGA package includes an array of external ball contacts, such as solder balls, that permit the package to be surface mounted to a printed circuit board (PCB) or other electronic component. Some BGA packages have a foot print that is about the size of the die contained in the package. These BGA packages are also known as chip scale packages.\nAs BGA packages are made smaller, and with higher lead counts, the electrical connections between the die and the ball contacts for the package become more difficult to make. The component of the package that makes the electrical connections between the die and package is sometimes referred to as an interconnect.\nOne type of prior art BGA package uses an interconnect in the form of a multi layered polymer tape, such as TAB tape. The polymer tape can be provided as a \"strip\" or \"coupon\" of material for packaging several dice at a time. Typically, the polymer tape comprises a polyimide substrate having patterns of metal conductors formed thereon. In addition, portions of the conductors can comprise beam leads, formed in a configuration which allows the beam leads to be bonded to the bond pads on the die.\nFIG. 1 illustrates a prior art BGA package 10. The BGA package 10 includes: a semiconductor die 12; a polymer tape 14 bonded to a face of the die 12; and an encapsulant 16 bonded to the face and sides of the die 12. The BGA package 10 also includes an adhesive layer 18 for bonding the polymer tape 14 to the die 12. In addition, the BGA package 10 includes a dense array of ball contacts 20 formed on the polymer tape 14. A solder mask 15 locates and electrically insulates the ball contacts 20 from one another.\nThe polymer tape 14 includes a pattern of metal conductors 21 that form separate electrical paths between bond pads 24 on the die 12 and the ball contacts 20. The conductors 36 include ball bonding pads 23 that are bonded to the ball contacts 20, and beam leads 22 that are bonded to the bond pads 24 on the die 12. The beam leads 22 are also encapsulated in the encapsulant 16.\nA representative process for forming the BGA package 10 includes the initial step of bonding one or more dice 12 to a strip of the polymer tape 14. The beam leads 22 can then be bonded to the device bond pads 24. Next, the encapsulant 16 can be formed, and the ball contacts 20 bonded to the ball bonding pads 23. The individual BGA packages 10 can then be singulated from the strip of polymer tape 14 and tested.\nTypically, a thermosonic bonding process, using gold, or gold plated materials, is employed to bond the beam leads 22 to the bond pads 24. Specialized bonding tools are required to make the bonds between the beam leads 22 and the bond pads 24. In addition, the beam leads 22 are subjected to stresses from the bonding and encapsulation processes, and during subsequent use of the BGA package 10. These stresses can cause the bonds to weaken and the beam leads to separate from the bond pads 24.\nBecause of these and other deficiencies in conventional BGA packages, improvements in BGA packages, and in methods for fabricating BGA packages, would be welcomed in the art. The present invention is directed to an improved interconnect for constructing BGA packages. The interconnect is constructed to facilitate formation of bonded connections with the die, and to provide improved bonded connections with the die."}
{"text": "Chronic myelogenous leukemia (CML) is a hematological stem cell disorder caused by increased and unregulated growth of myeloid cells in the bone marrow, and the accumulation of excessive white blood cells. Abelson tyrosine kinase (ABL) is a non-receptor tyrosine kinase involved in cell growth and proliferation and is usually under tight control. However, 95% of CML patients have the ABL gene from chromosome 9 fused with the breakpoint cluster (BCR) gene from chromosome 22, resulting in a short chromosome known as the Philadelphia chromosome. This Philadelphia chromosome is responsible for the production of the BCR-ABL fusion protein, a constitutively active tyrosine kinase that causes uncontrolled cellular proliferation. An ABL inhibitor, imatinib, was approved by the FDA for the treatment of CML, and is currently used as first-line therapy. It has been reported that 80% of CML patients respond to imatinib with under 3% progressing to advanced disease within 5 years (Schwartz, P. A., et al., Bioorg Chem 39: 192 (2011)). The durability of clinical response, however, is adversely affected by the development of resistance to drug therapy. In 2001, the first imatinib resistant mutant was reported as a T315I BCR-ABL “gatekeeper” mutation. Subsequent analysis revealed that reoccurrence arises with over fifty-five documented point mutations occurring throughout the catalytic and regulatory domains, with a large percentage located in the G-loop and the gatekeeper position (Id.). Clinical success has been achieved towards most mutations, but there is no approved agent directed towards the gatekeeper T315I mutation (O'Hare, T., et al., Clin Cancer Res 17: 212 (2011)).\nThe compound, N-methyl-2-[3-((E)-2-pyridin-2-yl-vinyl)-1H-indazol-6-ylsulfanyl]-benzamide or 6-[2-(methylcarbamoyl)phenylsulfanyl]-3-E-[2-(pyridin-2-yl)ethenyl]indazole, of the following structure:\nis known as axitinib or AG-013736.\nAxitinib is a potent and selective inhibitor of vascular endothelial growth factor (VEGF) receptors 1, 2 and 3. These receptors are implicated in pathologic angiogenesis, tumor growth, and metastatic progression of cancer. Axitinib has been shown to potently inhibit VEGF-mediated endothelial cell proliferation and survival (Hu-Lowe, D. D., et al., Clin Cancer Res 14: 7272-7283 (2008); Solowiej, S., et al., Biochemistry 48: 7019-31 (2009)). Clinical trials are currently on-going to study the use of axitinib for the treatment of various cancers, including liver cancer, melanoma, mesothelioma, non-small cell lung cancer, prostate cancer, renal cell carcinoma, soft tissue sarcomas and solid tumors. Inlyta® (axitinib) has been approved in the United States, Europe, Japan and other jurisdictions for the treatment of renal cell carcinoma.\nAxitinib, as well as pharmaceutically acceptable salts thereof, is described in U.S. Pat. No. 6,534,524. Methods of making axitinib are described in U.S. Pat. Nos. 6,884,890 and 7,232,910, in U.S. Publication Nos. 2006-0091067 and 2007-0203196 and in International Publication No. WO 2006/048745. Dosage forms of axitinib are described in U.S. Publication No. 2004-0224988. Polymorphic forms and pharmaceutical compositions of axitinib are also described in U.S. Publication Nos. 2006-0094763, 2008-0274192 and 2010-0179329. The patents and patent applications listed above are incorporated herein by reference."}
{"text": "1. Field of the Invention\nThe present invention relates generally to synthesizers and, more particularly, to direct digital synthesizers.\n2. Description of the Related Art\nA direct digital synthesizer generates analog signals that correspond to input tuning words and a synthesizer clock which also determines the synthesizer's Nyquist range. Although direct digital synthesizers provide a number of advantageous features (e.g., phase-continuous switching, high frequency resolution and fast switching time), they typically generate spurious signals which unduly limit their spurious free dynamic range. This limitation continues to be the most serious deficiency of direct digital synthesizers."}
{"text": "This invention relates to an improved means and method for connecting valves and the like to adjacent conduits, and more particularly to such improved means and methods utilizing a separate bolting ring threaded onto an end hub of the valve for connecting the valve to the conduit.\nHeretofore, such as illustrated in FIG. 1A herein, a separate bolting ring has been utilized for connecting an end hub of a valve to a flow line or conduit. However, with such an arrangement, it has been difficult to obtain precise and accurate positioning of the ring onto the hub. Further, upon relative high torques exerted by higher excessive bolting forces against the bolting ring, the thread area adjacent the flow line is placed under an undesirable compressive loading.\nOne possible solution to the precise positioning of a bolting ring on the externally threaded hub has involved the use of interference threads on the hub and bolting ring. However, this has tended to be unsatisfactory as it is difficult to install a bolting ring with an interference thread, and if the bolting ring is heated and then inaccurately positioned on the hub, the ring is fixed at that location unless removed by difficult procedures.\nAlso, if the threaded connection is slightly loose after heating of the bolting ring, and subsequent cooling, the bolting ring may rotationally deform under subsequent excessive bolting torques."}
{"text": "1. Field of the Invention\nThe invention relates to a vehicular rearview/sideview mirror assembly.\n2. Description of the Related Art\nRearview/sideview mirror assemblies are ubiquitous for contemporary vehicles, providing an occupant with a view from the vehicle of rearward and sideward areas collectively extending from the lateral limit of an occupant's peripheral vision to an exterior side of a vehicle, referred to generally as the “operational field of view.” Such mirror assemblies are typically mounted to the exterior of the vehicle at or proximate the front of the driver's and front passenger's doors (i.e. “side-mounted”). A conventional mounting location is adjacent the forward corner of the side windows. It is known to use afocal mirrors as rearview/sideview mirrors.\nTo adjust and expand the area of the rearward and sideward views observable by an occupant, the reflective element of a side-mounted mirror assembly can be configured for permanent or intermittent extension away from the vehicle exterior. However, this can project the reflective element into the vehicle airflow, resulting in drag, wind noise, and excess fuel consumption.\nConventional rearview/sideview mirror assemblies are typically unable to provide an occupant with an uninterrupted view of the entire area to the rear and side of a vehicle without moving the reflective element, or utilizing a plurality of reflective elements, each associated with a different field of view. Thus, a portion of the zone to the rear and side of a vehicle, referred to herein as the “hidden zone,” is not readily visible in the reflective element. Moving the reflective element is generally impracticable; utilizing a plurality of reflective elements will increase the aerodynamic drag force associated with a side-mounted mirror assembly.\nConventional side-mounted mirror assemblies can also incorporate several different elements, such as reflective elements, lenses, and the like, in order to maximize the operational field of view, while adhering to automotive safety and performance standards. However, the greater the number of elements making up a mirror assembly, the larger the mirror assembly housing must be to accommodate the elements, and the greater the aerodynamic profile associated with the projection of the mirror assembly from the vehicle exterior.\nIt would be desirable to utilize a rearview/sideview vehicle mirror assembly incorporating a minimum number of operational elements, and having a reduced weight and aerodynamic profile, capable of providing the operator of the vehicle with an enhanced view of objects, such as approaching vehicles, to the rear and side of the vehicle, which complies with automotive safety and performance standards."}
{"text": "The present invention relates to a combination pneumatic chuck having a position detecting mechanism for detecting a chuck position of a workpiece by jaw members.\nThere is a known combination pneumatic chuck having a chuck function of chucking a workpiece and a stroke function of moving the chucked workpiece forward and rearward as disclosed in Japanese Patent Application Laid-open No. 8-19983, for example. In this pneumatic chuck, a stroke air cylinder mechanism formed of a larger-diameter piston and a larger-diameter rod and a chuck air cylinder mechanism formed of a smaller-diameter piston and a smaller-diameter rod are coaxially incorporated in a double layers in inner and outer positions in a casing, a chuck head is mounted to the rod of the stroke air cylinder mechanism, a plurality of jaw members for grasping the workpiece are mounted to the chuck head such that the jaw members can be opened and closed, and the jaw members are opened and closed by the rod of the chuck air cylinder mechanism.\nTo such a pneumatic chuck, in general, a position detecting mechanism formed of a magnet and a magnetometric sensor is attached so as to detect chucking or releasing of the workpiece by the jaw members and to produce control signals for automatic operation. Specifically, magnets are mounted to the jaw members and a sensor is mounted to the casing to directly detect operating positions of the jaw members. Otherwise, a magnet is mounted to the piston of the chuck air cylinder mechanism to detect an operating position of the piston, thereby indirectly detecting opening and closing positions of the jaw members.\nHowever, because the chuck head for retaining the jaw members is driven by the stroke air cylinder mechanism in the above prior-art combination pneumatic chuck, it is difficult to directly detect the operating positions of the jaw members by mounting the sensor to the chuck head that is displaced. Because the piston and the rod of the chuck air cylinder mechanism are incorporated into the piston and the rod of the stroke air cylinder mechanism such that the stroke air cylinder mechanism exists between the chuck air cylinder mechanism and the casing, it is also difficult to detect the position of the piston of the chuck air cylinder mechanism.\nIt is a technical object of the present invention to provide a combination pneumatic chuck with a chuck air cylinder mechanism and a stroke air cylinder mechanism coaxially incorporated in inner and outer positions in the casing, wherein chuck operation of a workpiece by jaw members can be detected reliably by a simple structure.\nTo achieve the above object, a combination pneumatic chuck according to the invention comprises: a casing having therein a cylinder bore; a stroke first air cylinder mechanism including a larger-diameter first piston housed for sliding in the cylinder bore in the casing and a larger-diameter first rod having a base end portion connected to the first piston and a tip end portion extending from of the casing; a chuck second air cylinder mechanism including a smaller-diameter second piston housed for sliding coaxially in the first piston and a smaller-diameter second rod having a tip end portion connected to the second piston and extending to an inside of the first rod and a base end portion extending to such a position as to project from the first piston and to face a position detecting portion of the casing; a chuck head mounted to the tip end of the first rod for moving forward and rearward in response to expansion and contraction of the first rod; a plurality of jaw members for grasping a workpiece and mounted for opening and closing to the chuck head; a converting mechanism disposed between the jaw members and the second rod in the chuck head for converting expanding and contracting operations of the second rod into opening and closing operations of the jaw members; and chuck position detecting means which is for detecting opening and closing positions of the jaw members and which includes a detected member mounted to the base end portion of the second rod and a position sensor mounted to a position detecting portion of the casing for detecting the detected member.\nAccording to the combination pneumatic chuck of the invention having the above structure, by a simple structure in which: the base end portion of the second rod of the second air cylinder mechanism projects from the first piston of the second air cylinder mechanism and extends to such a position as to face the position detecting portion of the casing; the detected member is mounted to the base end portion; and the position sensor is mounted to the position detecting portion of the casing, the opening and closing positions of the jaw members can be indirectly detected by detecting the operating position of the rod of the first air cylinder mechanism even in the structure in which the first air cylinder mechanism is incorporated into the second air cylinder mechanism without extensively modifying the structure.\nAccording to a concrete embodiment of the invention, the position detecting portion is in a small cylindrical shape of such a size that the base end portion of the second rod is fitted into the position detecting portion and is formed to project outward along an axis at a central portion of the base end face of the casing, at least one sensor mounting groove is formed in an outer face of the position detecting portion, and the position sensor is mounted in the sensor mounting groove.\nAccording to another concrete embodiment of the invention, the pneumatic chuck includes stroke position detecting means which is for detecting a stroke position of the chuck head and which has a detected member mounted to the first piston and a position sensor mounted to the casing for detecting the detected member.\nAccording to a preferred embodiment of the invention, the pneumatic chuck includes a pair of stroke ports for supplying compressed air to pressure chambers on opposite sides of the first piston and a pair of chuck ports for supplying compressed air to pressure chambers on opposite sides of the second piston and the ports are concentrated on the base end face of the casing."}
{"text": "The last decade has seen knowledge of the immune system and its regulation expand tremendously. One area of particular interest has been that of research on the proteins and glycoproteins which regulate the immune system. One of the best known families of these molecules are the cytokines. These are molecules which are involved in the \"communication\" of cells with each other. The individual members of the cytokine family have been found to be involved in a wide variety of pathological conditions, such as cancer and allergies. Whereas sometimes the cytokines are involved in the pathology of the condition, they are also known as being therapeutically useful.\nInterleukins are one type of cytokine. The literature on interleukins is vast. An exemplary, but by no means exhaustive listing of the patents in this area includes U.S. Pat. No. 4,778,879 to Mertelsmann et al.; U.S. Pat. No. 4,490,289 to Stern; U.S. Pat. No. 4,518,584 to Market al.; and U.S. Pat. No. 4,851,512 to Miyaji et al., all of which involve interleukin-2 or \"IL-2.\" Additional patents have issued which relate to interleukin-1 (\"IL-1\"), such as U.S. Pat. No. 4,808,611 to Cosman. The disclosure of all of these patents are incorporated by reference herein. More recent patents on different interleukins include U.S. Pat. No. 5,694,234 (IL-13); U.S. Pat. No. 5,650,492 (IL-12); U.S. Pat. Nos. 5,700,664, 5,371,193 and U.S. Pat. No. 5,215,895 (IL-11); U.S. Pat. Nos. 5,728,377, 5,710,251, 5,328,989 (IL-10); U.S. Pat. Nos. 5,580,753, 5,587,302, 5,157,112, 5,208,218 (IL-9); U.S. Pat. Nos. 5,194,375, 4,965,195 (IL-7); U.S. Pat. Nos. 5,723,120, 5,178,856 (IL-6), and U.S. Pat. No. 5,017,691 (IL-4). Even a cursory review of this patent literature shows the diversity of the properties of the members of the interleukin family. One can assume that the larger cytokine family shows even more diversity. See, e.g., Aggarwal et al., ed., Human Cytokines: Handbook For Basic And Clinical Research (Blackwell Scientific Publications, 1992), Paul, ed., Fundamental Immunology (Raven Press, 1993), pg 763-836, \"T-Cell Derived Cytokines And Their Receptors\", and \"Proinflammatory Cytokines and Immunity.\" All cited references are incorporated by reference.\nThe relationships between various cytokines are complex. As will be seen from the references cited herein, as the level of a particular cytokine increases or decreases, this can affect the levels of other molecules produced by a subject, either directly or indirectly. Among the affected molecules are other cytokines.\nThe lymphokine IL-9, previously referred to as \"P40,\" is a T-cell derived molecule which was originally identified as a factor which sustained permanent antigen independent growth of T4 cell lines. See, e.g., Uyttenhove et al., Proc. Natl. Acad. Sci. 85: 6934 (1988), and Van Snick et al., J. Exp. Med. 169: 363 (1989), the disclosures of which are incorporated by reference, as is that of Simpson et al., Eur. J. Biochem. 183: 715 (1989).\nThe activity of IL-9 was at first observed on restricted T4 cell lines, failing to show activity on CTLs or freshly isolated T cells. See, e.g., Uyttenhove et al., supra, and Schmitt et al., Eur. J. Immunol. 19: 2167 (1989). This range of activity was expanded when experiments showed that IL-9 and the molecule referred to as T cell growth Factor III (\"TCGF III\") are identical to MEA (Mast Cell Growth Enhancing Activity), a factor which potentiates the proliferative response of bone marrow derived mast cells to IL-3, as is described by Hultner et al., Eur. J. Immunol. and in U.S. patent application Ser. No. 498,182 filed Mar. 23, 1990, the disclosures of both being incorporated by reference herein. It was also found that the human form of IL-9 stimulates proliferation of megakaryoblastic leukemia. See Yang et al., Blood 74: 1880 (1989). Recent work on IL-9 has shown that it also supports erythroid colony formation (Donahue et al., Blood 75(12): 2271-2275 (Jun. 15, 1990)); promotes the proliferation of myeloid erythroid burst formation (Williams et al., Blood 76: 306-311 (Sep. 1, 1990); and supports clonal maturation of BFU-E's of adult and fetal origin (Holbrook et al., Blood 77(10): 2129-2134 (May 15, 1991)). Expression of IL-9 has also been implicated in Hodgkins's disease and large cell anaplastic lymphoma (Merz et al., Blood 78(8): 1311-1317 (Sep. 1, 1990). Genetic analyses of mice that were susceptible or resistant to the development of bronchial hyperresponsiveness have unraveled a linkage with the IL-9 gene as well as a correlation between IL-9 production and susceptibility in this model (Nicolaides et al., Proc. Natl. Acad. Sci. USA, 94, 13175-13180, 1997). Human genetic studies also point to the IL-9 and IL-9R genes as candidates for asthma (Doull et al., Am. J. Respir. Crit. Care Med., 153, 1280-1284, 1996; Holroyd et al., Genomics 52, 233-235, 1998). Secondly, IL-9 transgenic mice allowed for the demonstration that increased IL-9 expression result in lung mastocytosis, hypereosinophilia, bronchial hyperresponsiveness and high levels of IgE (Temann et al., J. Exp. Med. 188, 1307-1320, 1998; Godfraind et al., J. Immunol. 160, 3989-3996, 1998; McLane et al., Am. J. Resp. Cell. Mol. 19:713-720 (1999). Taken together, these observations strongly suggest that IL-9 plays a major role in this disease Additional work has implicated IL-9 and muteins of this cytokine in asthma and allergies. See, e.g. PCT Application US96/12757 (Levitt, et al), and PCT US97/21992 (Levitt, et al), both of which are incorporated by reference..\nIL-9 is known to affect the levels of other molecules in subjects. See Louahed et al., J. Immunol. 154: 5061-5070 (1995; Demoulin et al., Mol. Cell. Biol. 16: 4710-4716 (1996), both incorporated by reference. It will be recognized that the molecules affected have their own functions in biological systems. For example, Demoulin et al. show that many of the known activities of IL-9 are mediated by activation of STAT transcription factors. As such, there is continued interest in trying to identify molecules whose presence and/or level is affected by other molecules, such as cytokines.\nThe disclosure which follows describes such molecules. It was found that nucleic acid molecules encoding the proteins of the invention were expressed in the presence of IL-9, but not in its absence. Hence, these molecules are, inter alia, \"markers\" for the expression or effect of IL-9 in a subject. The molecules are referred to as T Cell Derived Inducible Factors or \"TIFs\" hereafter. These and other features of the invention will be seen in the disclosure which follows."}
{"text": "1. Field of the Invention\nThe present invention relates to a folding device and, in particular, to a folding device able to be mounted on a vehicle for folding and/or locking front and rear frames of the vehicle.\n2. Description of the Related Art\nIn order to save storage space and facilitate carrying, the current trend is to build folding vehicles, such as folding bicycles or electric folding bicycles, and in order to meet different requirements in respect of folding, various folding devices have been designed on the foldable vehicles, among which folding, unfolding and locking devices are the most commonly used.\nFurther, the conventional folding device substantially includes two linking members and a quick release skewer. The two linking members are respectively jointed to front and rear frames of the folding vehicle by welded or other connection means, and the two linking members are pivotally connected with each other by means of a suitable mechanism. The quick release skewer includes a rod threaded on one end and with a lever operated cam assembly on the other. The rod is inserted into the two linking members, a nut is threaded on, and the lever can be closed to selectively tighten the cam and secure the two linking members to each other.\nHowever, when the two linking member are intend to pivot with each other for folding, a user must open and rotate the lever and then pull the rod out from the two linking members to loosen the cam. The operation of the conventional folding device is monstrously complicated, and the connection strength between the two linking member is not strong enough.\nThe present invention is, therefore, intended to obviate or at least alleviate the problems encountered in the prior art."}
{"text": "The present invention relates to a pixel structure for use in an infrared pixel array. More specifically, the structure of the present invention is used for microemitter or microbolometer applications. In the microemitter application, the pixel of the present invention is used as one portion of a pixel array to perform infrared projection. Alternatively, in the microsensor, or microbolometer applications, the pixel structure of the present invention forms one pixel of a detector array.\nMicrostructures generally can be used to form various types of sensors and sensor products. In certain applications, these microstructures are arranged in arrays which cause the multiple pixels to cooperate with one another. One well known application is the microemitter, where the pixel array is used to project an image. In another well known application, the pixel array is used as a microbolometer wherein the array is used to detect a two dimensional signal.\nIn each application, it is essential that the pixels be isolated from one another. This isolation allows each pixel to operate independently from its surroundings pixels, thus creating contrast in either the sensed or displayed image.\nIn the aforementioned microemitter application, each pixel is independently energized to project a signal. In the case of an IR scene projector, the pixel is energized, causing heating of the element itself resulting in the emission of IR radiation. When arranged in an array of individual pixels and appropriately addressed, this collection of individual emitters can create the desired image.\nSimilarly, an array of microstructure pixels can be used in the microbolometer application to detect IR radiation or signals. In this application, each individual pixel is sensitive to IR signals. When such IR signals hit the sensing surface or area of the pixel, the resistance of the pixel active element changes. This resistance change can be appropriately sensed by related circuitry and fed back to an image controller. When the signals from an array of pixels are appropriately processed, a digital picture image can be created.\nFor any of the microstructure devices referenced above (microbolometer or microemitter), the optimum design is influenced by many factors. Due to the practical operational requirements, these considerations are surprisingly complex. More specifically, the optimum design characteristics include large fill factors, high thermal isolation, and an appropriate thermal time constant (speed).\nGenerally speaking, the time constant is related to the speed of the device. As can be appreciated, in either a detection or projection mode, it is desired that the pixels are updated frequently. This requires that the pixels be able to cycle or regenerate very quickly. In practical applications, it is not unusual to have frame rates existing anywhere from 30 to 500 hertz, thus the appropriate time constant must exist. Generally speaking, the thermal time constant is equal to the thermal mass of the pixel divided by the thermal conductivity. A low thermal time is desired as this will allow for high speed operation.\nAs mentioned above, large fill factor is an optimum characteristic in the design of microstructure pixels. Fill factor relates to the area used by the active portion of the pixel. It is beneficial for the fill factor to be as high as possible, utilizing the largest potential area of the pixel. This is especially true for emitter applications, where it is desirable to have a large surface of the pixel emitting the desired radiation signals.\nAnother obvious consideration is thermal isolation. It is clearly necessary that each of the pixels be thermally isolated from one another in order to avoid any cross talk. This allows each pixel to maintain its independence and create a high contrast array.\nThe ability to distribute and appropriately transmit heat is of the utmost importance in designing the pixel. The ability to dissipate heat when the heated pixel is turned off clearly affects the thermal time constant and speed of the pixel. When designing the microstructure itself, heat dissipation considerations must always be taken into account.\nCurrent microstructures for IR pixels generally can take on several configurations. Each of these configurations however has several consistent features, all related to the manufacturing process. As can be expected, the microstructures used for these infrared pixels are fabricated using thin film processes. Consequently, appropriate films and masks are used in a multi-step process to achieve the appropriate microstructure configuration.\nOne example of a manufacturing process for creation of these microstructures starts with integrated read out and drive circuitry, previously fabricated on a wafer. This wafer is then finished to have a planar top surface. Next a reflective layer is placed upon this finished top surface. Upon the reflective layer is deposited a sacrificial layer with connection posts at appropriate points to allow communication with the integrated circuit. Upon the sacrificial layer is deposited the actual microstructure itself which includes appropriate layers of resisted material and silicon nitride. When the sacrificial layer is then removed, the microstructure is left, appropriately coupled to the control electronics via connective posts and positioned slightly above the reflective layer. Generally speaking, the active emitter or detector material is encapsulated within an insulating film, such as silicon nitride. This coating provides for mechanical support and chemical passivation. The pixels have a uniform height, generally equal to the height of the silicon nitride layer.\nThe above-referenced manufacturing process creates pixels of adequate operation, however, their design is not optimized. It is desired to provide the ability to increase the speed at which the pixels can operate without compromising other performance criteria."}
{"text": "As a so-called writable optical recording medium, there has been known an optical recording medium which comprises a transparent substrate, a light absorbing layer comprising an organic coloring matter (or colorant) such as a cyanine type and a phthalocyanine type, a light reflection layer formed on the light absorbing layer, and a protection layer disposed on the light reflection layer for preventing the deterioration or scratch which can occur in the light absorbing layer and/or the light reflection layer. In the optical recording medium having such a structure, a laser light (or laser beam) is supplied thereto from the transparent substrate side thereof and the laser light is concentrated (or focused) on a minute area of the light absorbing layer so that it is converted into heat energy. As a result, the thus irradiated portion of the light absorbing layer is caused to have a different state or configuration (e.g., a pit is formed in such a portion), thereby to effect writing. On the other hand, a laser light (or laser beam) for reading is supplied to such an optical recording medium from the transparent substrate side thereof and is reflected by the light reflection layer. At this time, the resultant contrast between the quantity of the reflected light in a portion having the above pit formed therein and the quantity of the reflected light in a portion having no pit formed therein is detected and converted into an electric signal, thereby to effect reading.\nIn the conventional optical recording medium, however, the above protection layer has been formed by applying a resin material such as an ultraviolet ray curing (or hardening) resin and a solvent type resin on the light reflection layer and hardening the resultant coating, and the resultant protection film only has a thickness of about 5 to 15 .mu.m. Accordingly, when such a protection layer is disposed on the light reflection layer by itself, the protection layer can easily peeled or separated from the light reflection layer when an external force such as scratching with a nail of a user. As a result, in the conventional optical recording medium, there is posed a problem such that the light absorbing layer and the light reflection layer are not sufficiently protected from the external force as described above. In addition, when the protection layer is caused to have a larger thickness, the above problem has not been solved fundamentally or drastically."}
{"text": "The invention concerns a device that is a cultivation tank for aquatic organisms.\nFish tanks are presently either made of plastic, steel, or concrete in a number of designs. The removal of excess feed, excrement and the like is usually done by emptying large volumes of water from the tank to allow technical staff to brush the bottom of the tank clean. If the water outlet and supply is forced to flow by means of a current, there is a better degree of self-cleaning. This however requires considerable amounts of water to be supplied.\nMarine fish larva such as cod, turbot, halibut and plaice are too sensitive for such strong water current. Thus the cleaning of tanks for these species has to be done manually by technical staff vacuuming the bottom with a syphon that has a suitable nozzle fitted to the end. This operation is excessively time consuming. In cases where the water is recirculated it is disadvantageous to have a strong water current flowing through the tank."}
{"text": "The present invention relates generally to apparatus and methods for the purpose of determining the progress of a food load being cooked or otherwise heated in a microwave oven. More particularly, the invention relates to such apparatus and methods responsive to changes in electrical characteristics of the food load as it is heated, particularly changes as a result of variations in state and quantity of moisture content.\nIt is desirable in a cooking microwave oven to be able to monitor or determine the progress of food being cooked or heated. While microwave ovens typically include viewing windows, such windows provide limited visibility, and often little information is obtained concerning the progress of food cooking by observing the appearance in any event. One general method, useful particularly in the case of large pieces of meat being cooked, is to employ a temperature-sensing probe assembly inserted into the food being cooked, such as is disclosed in the Chen et al. U.S. Pat. No. 3,975,720 and in the Meek et al. U.S. Pat. No. 4,086,813. However, the use of such a probe is not always convenient or possible, and it desirable to provide further alternative approaches. (As employed herein, the term \"cooking\" is employed in a broad sense to mean the heating or thermalization of food placed in a microwave oven, regardless of the particular temperature range over which the heating occurs and regardless of the particular chemical or physical change occurring within the food.)\nAnother general method of providing information about the progress of food cooking in a microwave oven relies upon changes in the characteristics of the food as an electrical load or a microwave absorber during cooking. For example, conductivity and dielectric properties of food change during cooking or heating, particularly as water content is affected by the microwave heating. One known approach relying upon such changes monitors the standing wave ratio or a related property within a feed waveguide or other form of transmission line supplying the microwave cooking cavity. This general approach is disclosed for example in the Moe U.S. Pat. No. 3,813,918.\nAlthough not directly responsive to changes in the load properties of individual items of food as cooking progresses, similar sensing principles have been proposed for controlling the length of cooking time as a function of energy delivered to a food load or available energy apportioned between a plurality of individual items of food. For example, in the system of the Schroeder U.S. Pat. No. 2,744,990, cooking time is related to net energy supplied to the microwave cooking cavity as determined by a directional coupler in a feed transmission line, net power delivered being sensed forward power minus sensed reflected power. Similarly, in the systems described in the Moore U.S. Pat. No. 3,999,027 and Tallmadge et al. U.S. Pat. No. 4,009,359, microwave field strength is sensed within the cooking cavity itself, rather than in a feed waveguide, for the purpose of controlling the time duration of operation to achieve a desired temperature within a food load material. With a higher microwave field strength, it is assumed that the cooking effect is greater, and the time duration is accordingly shortened. Specifically, electromagnetic field strength is integrated with respect to time as an indicator of overall cooking effect.\nAnother condition which may be sensed using related techniques, particularly for protective purposes, is the absence of any food load whatsoever within the microwave cooking cavity. Under such conditions, the standing wave ratio within the feed waveguide, as well as the field strength within the cavity, are higher than normal. Various systems have been proposed for sensing such conditions, and automatically turning off the microwave generator in response. For example, the system of the Meissner et al. U.S. Pat. No. 3,281,567 directly senses field strength within a cooking cavity. In a somewhat similar fashion, the system of the Haagensen et al. U.S. Pat. No. 3,527,915 indirectly responds to field strength within a cooking cavity by sensing the temperature rise of an element placed at the bottom of the cavity and which absorbs microwave energy. Other protective systems respond to conditions within the feed waveguide, for example, as disclosed in the Anderson U.S. Pat. No. 4,412,227, the Kohler et al. U.S. Pat. No. 3,491,222, the Jones et al. U.S. Pat. No. 3,662,140, and the Bucksbaum U.S. Pat. No. 3,670,134.\nIn a somewhat different vein, a related principle of operation is utilized in an automatic control system for a continuously-moving type microwave dryer disclosed in the Kashyap et al. U.S. Pat. No. 4,035,599. In the Kashyap et al. system, microwave energy is passed through a moving web load, such as paper. Input power and output power are separately sensed by means of directional couplers. In order to control microwave input power to maintain a constant drying effect, microwave input power is varied as a function of sensed input and output power, as well as of web velocity.\nThe effect of varying quantities of food and changes within a food load as cooking progresses is recognized in the systems of the Sawada U.S. Pat. No. 3,104,304 and the Stecca et al. U.S. Pat. No. 3,321,604, each attempting to maintain optimum conditions as the food changes. In Sawada, the oscillator frequency is varied to maintain resonance. In Stecca et al., the cavity itself is tuned to maintain resonance.\nVarious non-electrical approaches to the problem of monitoring cooking progress in a microwave oven have also been proposed. For example, the Smith U.S. Pat. No. 3,467,804 proposes a sensor for steam or other vapors which may be emitted when an article of food is heated, or has reached a predetermined temperature. In the Ueno U.S. Pat. No. 4,049,938, an infrared radiation detector is proposed.\nWhile the various approaches described above for indirectly obtaining information concerning the progress of food being cooked in a microwave oven, particularly those which rely upon a change in the load properties of the food itself, do function to some extent, greater sensitivity is desirable. This greater sensitivity is provided by the apparatus and method of the present invention. In particular, the sensing effected by the present invention provides a relatively large percentage change in the sensed parameter depending upon the dielectric properties of the food load, particularly as a result of the state and quantity of its moisture content."}
{"text": "Air conditioners are provided on their front panel with an air duct, where is attached an arrangement of air flow deflecting vanes, which allow the air flow to be directed according to predetermined directions.\nIn a known prior art construction, the arrangement of air flow deflecting vanes to be provided in the air duct of the front panel of an air conditioner comprises a plurality of deflecting vanes, each provided with a pair of pivoting pins, which are oppositely disposed to each other and aligned so as to coincide with a rotating axis of the vane, said pins individually engaging each deflecting vane to the front panel.\nIn this construction, the holes made in the air outlet of the front panel and which receive the pivoting pins are in the form of cuts made in a wall portion which defines the contour of the air outlet opening of the front panel during the manufacturing process of the latter, usually by injection. This prior art construction has operational disadvantages, such as each vane being individually manually mounted to the front panel and the wear of said holes caused by use, said wear resulting in a diametrical variation of the holes, enlarging the diameter thereof and allowing gaps to arise. These gaps, which may exist since the molding process, cause vibration of said vanes, provoking noise. This vibration problem is particularly intense in this construction, due to the vanes being directly applied to the front panel by means of a pin-hole fitting. in another known solution, each vane of the vane arrangement is produced incorporated to a single driving bar by corresponding weakened hinge portions. The assembly of the vane arrangement and driving bar to the front panel is carried out by fitting a pivoting pin provided from a lateral edge of each corresponding deflecting vane, opposite to the other lateral edge of the latter connected to the bar by the hinge portion, each said pin being fittable in a corresponding hole produced in the front panel and coinciding with a rotating axis of the respective deflecting vane.\nThis construction has the disadvantage of producing noise during operation and use, caused by vibrations resulting from the gaps which appear in the pin-hole fitting, due to the vanes being directly applied to the front panel by a pin-hole assembly.\nIn these constructions, it is common to occur variations in the frame of the front panel when the latter is affixed to the structure of the air conditioner, which lead to a disalignment between the upper and lower edges of said front panel. Said disalignment creates gaps which may lead to the disengagement of the vanes fitted in said panel and to the release of said vanes.\nFurther to these disadvantages, the known prior art constructions have the inconveniences of requiring a large number of perforations in the front panel, corresponding to the amount of hinge pins for each vane to be affixed to said front panel, besides the fact that the vane arrangement is defined by each vane being individually mounted to the front panel, turning the process slow and costly."}
{"text": "(1) Field of the Invention\nThe present invention relates to a scheduler, a radio base station apparatus including the scheduler, and a scheduling method. The invention relates to, for example, technology suitable for use in a radio base station apparatus in a radio mobile communication system.\n(2) Description of the Related Art\nIn a radio mobile communication system using packet communication, communication quality is monitored, and control is performed so that the transmission speed and the modulation method suitable for communication quality are adaptively selected. The reason of change in communication quality is variation in radio propagation environment, and it is often occurred that radio quality is instantly becomes lower than the required radio quality in the mobile radio system.\nWhen radio propagation environment is deteriorated, an error rate in a radio segment becomes high. Thus, in the radio layer, methods of lowering the error rate are performed such as decreasing transfer information amount, increasing the error correction information, or using a modulation scheme high in error resilience. In addition, when the error rate in a radio segment becomes high, opportunities of transmission of NACK which indicates abnormal reception on the reception end are increased, and the turn-around of the packet is lengthened. Therefore, in the upper layer, measurement of ACK/NACK and measurement of turn around time of a packet are performed. In this manner, communication quality is monitored, thereby performing controlling of allocation of radio resources (sub-carrier frequency, etc.) to satisfy the required quality.\nFurther, in MAC (Media Access Control) layer or a physical layer, using a scheduler which performs allocation of radio resources, the priority in consideration of radio quality is calculated. In accordance with the priority, controlling to allocate radio resources to users is performed, so that the system use efficiency (use efficiency of radio resources) is improved. Hereinafter, a function for calculating the priority will be called an evaluation formula.\nAs scheduling for allocating radio resources, there is a representative algorism of Maximum Carrier-to-Interference-Ratio (MaxCIR) (see paragraph [0006] of the following patent document 1). The MaxCIR method is a method in which radio resources are always allocated to a user who is the best in instantaneous radio quality [CIR or SINR (Signal-to-Interference-plus-Noise-power-Ratio) at the time scheduling is performed. According to the MaxCIR method, since radio resources are allocated to users good in radio quality, the error rate in a radio segment is low and the system throughput becomes the maximum. As to the use efficiency of the system, the MaxCIR method realizes the highest system throughput. Hence, it is the method realizing the outstanding merit for an operator of the radio system.\nHowever, the MaxCIR method causes a blind sector in which communication is not available, because radio resources are not allocated to users located in an area, such as a cell edge, in which radio quality is poor. Hence, in a case where plane-like service such as a cellular system is required, additional base stations for compensating for the blind sector are required.\nFurther, the Proportional Fairness (PF method) is also a well-known scheduling algorism (see paragraph [0007] of patent document 1 and the following non-patent document 1). According to the PF method, radio resources are allocated to users with good instantaneous radio quality with respect to average radio quality. Thus, since radio resources are allocated to all the users in a similar manner regardless of the instantaneous radio quality, fairness among the users is high. Hence, since radio resources are allocated to users located at cell edges, the number of blind sector becomes small, so that plane-like service is realized with base stations lower in number than MaxCIR.\nFurther, an algorism called the Generalized PF (GPF) method (see the following non-patent document 2) is also well-known. According to the GPF method, using the following formula (1), radio resources are allocated to a user having the highest result of the evaluation formula.\n                    value        =                              S            ins            α                                S            ave            β                                              (        1        )            \nHere, Sins indicates instantaneous radio quality; Save indicates average radio quality. The case of α=1, β=0 corresponds to the MaxCIR method, and the case of α=β=1 corresponds to the PF method. That is, the GPF method is a method in which balance between improvement and fairness of the system user efficiency can be changed by means of parameters α and β on the basis of the PF method. For example, if β is greater than a, users with lower average radio quality are allocated at high probability, so that even users whose quality is poor are allocated with high probability, and fairness of throughput among users is improved.\nHere, when such fairness is assured, the priority of users poorer in radio quality are higher, and thus, the fairness is assured with low throughput, so that the system use efficiency becomes poor.\n[Patent Document 1] Japanese Patent Application Laid-open No. 2003-152630\n[Non-patent Document 1] A. Jalai, R. Padovani, R. Pankaj, “Data Throughput of CDMA-HDR a High Efficiency-High Data Rate Personal Communication Wireless System”, VTC2000 Spring, May 2000\n[Non-patent Document 2] Nortel Networks, “Nortel Networks' reference simulation methodology for the performance evaluation of OFDM/WCDMA in UTRAN”, R1-03-0785, 3GPP TSG RAN WG1#33, August 2003\nTo use radio resources effectively, although it is important to monitor communication quality so as to control transmission using appropriate transmission speed and modulation method, instantaneous radio quality exceeding the predetermined error rate can appear due to instantaneous variation. According to the previous methods, even if the instantaneous variation becomes lower than the required radio quality, radio resources are allocated when the priority (hereinafter will be called the evaluation value) of radio resource allocation calculated by the scheduler is high. Even in such a case, since the upper layer monitors and controls communication quality, the radio quality becomes stable in average. However, since the upper layer is in charge of controlling, the control delay is large. When the control delay is large, instantaneous variation of radio quality cannot be followed up. Thus, radio resources are used in vain during control delay. In addition, since NACK occurs a lot, the traffic is suppressed.\nOn the other hand, to realize plane-like service, a scheduling method with a coverage as large as possible needs to be employed. Although in the PF method, the system throughput is lowered since radio resources are allocated to users low in radio quality, an operator is forced to use the PF method in order to ensure the coverage of a cell."}
{"text": "This invention relates to a vacuum coating system of the in-line load-lock type wherein substrates, such as silicon wafers for example, are transported for positioning within a vacuum chamber in proximity to a vapor source, such as sputtering electrodes or electron-beam evaporators, for coating as required.\nThe substrates which are first positioned on a carrier support at a loading port outside the vacuum chamber, are transported by the carrier through controlled gated load-locks and buffer zones into the process or coating chamber so as to be properly positioned with respect to, or slowly passed by the coating source. A vacuum coating system of this general type is disclosed in U.S. Pat. No. 3,662,708 granted May 16, 1972 to Robert L. Shrader and assigned to the same assignee as the present invention.\nIn this type of vacuum coating system, it is commercially advantageous to transport and position the workpieces or substrates without undue loss of production time; that is, it is desirable that the substrate transfer time be as small as possible in relation to the required coating time and that the equipment therefor be simple, positive and reliable. In many prior systems where the carrier is manually transferred from one position to another, or where one or more sets of drive rollers moving at constant speeds are used, or where complicated mechanical positioning means such as pneumatic thrust pistons, cables, linkages, etc. are used, the production process in general has not been satisfactory. This is especially the case where significant maintenance and down-time are encountered in complicated equipment, and where transfer of the substrates between process steps is time-consuming.\nThe present invention, therefore, is concerned with improved and simplified workpiece transfer equipment for a vacuum processing system of the character described, that readily lends itself to efficient transfer of the workpieces between load-locks, buffer zones, process chamber, etc. throughout the entire vacuum processing system."}
{"text": "Face gears are ring gears with a face angle (and root angle) equal to the shaft angle between the face gear and it's mating member. The mating member is a regular cylindrical spur or helical pinion. Standard face gears have a face angle of 90° which corresponds to a shaft angle of also 90° (such 90° gears are also known as “crown” gears).\nTo date, methods of manufacturing face gears have been complicated, with special tools dedicated to a single design applied on machine tools which are usually modified cylindrical gear manufacturing machines. Such methods for the soft manufacturing of face gears include:                Hobbing, using a job dedicated special hob on a cylindrical hobbing machine, which is modified in order to allow for cutting at the lowest circumferential section of the hobbing tool (vertical hobbing machine table axis).        Shaping, using a shaper cutter representing the mating cylindrical pinion, and a shaping machine, with a work table which is rotated (versus a regular cylindrical gear shaping machine) by the face gear set's root angle (commonly 90°).        Universal milling method, using an end mill on a 5-axes machining center.        Grinding from solid, using the grinding methods mentioned in the hard finishing section below.        \nToday's known methods for the hard finishing of face gears include:                Continuous grinding, using a threaded grinding wheel with a thread reference profile, identical to the face gear set's pinion tooth profile on a large diameter wheel with small width having generally 1.5 to 2.5 thread revolutions (see WO 98/02268; U.S. Pat. No. 6,390,894; and U.S. Pat. No. 6,951,501 the disclosures of which are hereby incorporated by reference).        Single index generating grinding with a wheel profile identical to the face gear set's pinion tooth profile.        Skiving, using a shaper cutter or special hob.        Skiving, using and end mill on a 5-axes machining center.        Honing, using a modified pinion with an abrasive layer on the tooth surface.        \nPresently, face gear soft machining methods depend on job specific, special tools, which are expensive and not flexible regarding their use for other jobs or for optimizations. The machining time of a face gear is in general significantly longer than the cutting time of a comparable cylindrical or bevel ring gear.\nTwo of the more common face gear hard finishing methods use either a very complex tool geometry which is difficult to dress and requires a long dressing time (threaded wheel grinding), or a complicated and time consuming generating roll, combined with a feed motion in face width direction (single index generating grinding).\nSkiving with a special hob or a shaper cutter made from carbide material provides reasonable cutting times but requires a tool which is not only expensive but also not readily available or not available at all.\nFace gear honing requires, for example, a heat treated, ground and CBN coated pinion, which is expensive, not flexible and depends on a rather large pinion offset (equal the required offset between face gear and mating cylindrical pinion) for good chip removal, which limits the application to face gear sets which have such a high offset."}
{"text": "Cationic d.sup.o metal complexes of the formula L.sub.2 MCH.sub.3.sup.+, where L is a cyclopentadienyl-type ligand and M is a Group IVB metal (titanium, zirconium, or hafnium) have been implicated as the active species in metallocene-based Ziegler-Natta olefin polymerization catalyst systems of the type L.sub.2 MCl.sub.2 /Al.sub.n Cl.sub.3-n.\nInterest in this type of catalyst resumed in the early 1980's and a variety of zirconocene complexes were found to form extremely active olefin polymerization catalysts when they were combined with methylaluminoxane (MAO). The activity increases with an increase of the Al/Zr ratio, while MAO itself cannot polymerize ethylene. The neutral compounds L.sub.2 MR (where L is as defined above, M is lutetium or scandium) and (L.sub.2 MH).sub.2 (where M is lanthanum, neodymium, samarium, or lutetium), which are isoelectronic with L.sub.2 MCH.sub.3.sup.+ (where M is titanium, zirconium, or hafnium), were also found to be highly active ethylene polymerization catalysts. On the other hand, the cationic titanocene complex could be trapped by the insertion of TMSC=CPh into the titanium-methyl bond. Unfortunately, such a cationic complex is catalytically inactive, possibly due to the bulkiness of the highly substituted vinyl ligand. Direct isolation of the cationic titanium methyl complex stabilized by conventional anions such as [BF.sub.4 ] or [PF.sub.6 ] can only be achieved from acetonitrile solutions resulting in the adducts of the formulae (L.sub.2 MCH.sub.3 (CH.sub.3 CN)).sup.+ (BF.sub.4) or (L.sub.2 MCH.sub.3 (CH.sub.3 CN)).sup.+ (PF.sub.6).sup.-, where M is zirconium or titanium and L is a cyclopentadienyl-type ligand. These are deactivated due to the ligand coordination of acetonitrile. Attempts to exchange acetonitrile with a weaker base, e.g., tetrahydrofuran (\"THF\"), produces L.sub.2 MFCH.sub.3 due to facile fluorine abstraction.\nBy using [BPh.sub.4 ].sup.- as the counterion, however, active cations can be obtained. The THF-coordinated cationic zirconocene complex of the formula (L.sub.2 ZrCH.sub.3 (THF)).sup.+ (B(C.sub.6 H.sub.5).sub.4).sup.- was isolated by the reaction of dicyclopentadienyl zirconium dimethyl with silver tetraphenylborate in acetonitrile followed by recrystallization from tetrahydrofuran. The NMR spectrum and the X-ray structure clearly show the coordination of the tetrahydrofuran. Tetrahydrofuran-coordinated cationic titanocene complexes of the formula (L.sub.2 TiCH.sub.3 (THF)).sup.+ (B(C.sub.6 H.sub.5).sub.4).sup.- were synthesized later. All of these complexes are weakly active ethylene polymerization catalysts (N.sub.t =7.2 min.sup.-1 and 1.9 min.sup.-1, respectively, at room temperature). It is worth noting that a similar reaction does not take place for analogous thorium complexes under mild conditions.\nWhile the concept of an active cationic center was established, research interest was gradually passing from confirming the cationic nature of the catalysts to searching for more active catalysts. The effect of base-coordination on the activity of the cationic catalysts has been studied systematically, and it is commonly believed that the coordination of the Lewis base decreases the catalytic activity of the cationic species. Extensive efforts to isolate base-free catalysts were undertaken. The pivotal concept of such research is to develop an anion which is bulky, highly stable, and noncoordinating to the cationic catalytic center, which exhibits very strong Lewis acidity and very high reactivity. Attempts at the isolation of the base-free cationic zirconocene complex, (L.sub.2 ZrCH.sub.3).sup.+ (B(C.sub.6 H.sub.5).sub.4).sup.-, failed and produced instead the neutral compound, L.sub.2 ZrCH.sub.3 (C.sub.6 H.sub.5), by the ligand redistribution reaction between the anion and the cation. The Lewis base-free cationic thorium complex (L.sub.2 ThCH.sub.3).sup.+ (B(C.sub.6 H.sub.5).sub.4).sup.- was synthesized as the first example of the base free catalysts by quantitative protonolysis of the neutral thorium alkyl complex L.sub.2 Th(CH.sub.3).sub.2 with ((C.sub.2 H.sub.5).sub.3 NH).sup.+ B(C.sub.6 H.sub.5).sub.4).sup.- in noncoordinating solvents. However, the low temperature solution NMR data and room temperature solid state NMR data show that the cation and the anion interact strongly with each other as a tight ionic pair. This could be the reason for the modest activity of the catalyst in ethylene polymerization (N.sub.t =1 min.sup.-1 at room temperature). Protonolysis of dicyclopentadienyl zirconium dimethyl with ((C.sub.4 H.sub.9).sub.3 NH).sup.+ (B(C.sub.6 H.sub.4 R).sub. 4).sup.-, where R can be hydrogen, methyl, or ethyl, gives a zwitterion of the formula L.sub.2 Zr.sup.+ (m-C.sub.6 H.sub.3 R)B(C.sub.6 H.sub.4 R).sub.3.sup.-. The activity of this catalyst for ethylene polymerization is 3.75.times.10.sup.5 gPe/mol.atm.h. (N.sub.t =3.75s.sup.-1) at 80.degree. C. The X-ray structure shows an agostic interaction between the cationic zirconium center and the C-H bond ortho to the boron.\nThe other type of noncoordinating anions that have been used are carborane-derived. Dicyclopentadienyl zirconium dimethyl reacts with the diprotic carborane complex nido-C.sub.6 B.sub.9 H.sub.13, giving off methane and resulting in (L.sub.2 ZrCH.sub.3).sup.+ (C.sub.2 B.sub.9 H.sub.12).sup.- as an active ethylene polymerization catalyst (N.sub.t =2.7s.sup.-1 at 40.degree. C.). The X-ray structure shows that the cationic zirconium center is bound to the anion through a terminal hydride on the nido face of the anion. Metallocarboranes of the formula (M(C.sub.2 B.sub.9 H.sub.11).sub.2).sup.-,where M is iron, cobalt, or nickel are also stable anions for the zirconocene and thorocene cations. For example, (L.sub.2 ZrCH.sub.3).sup.+ (M(C.sub.2 B.sub.9 H.sub.11).sub.2).sup.-, where M can be one of the aforementioned metals, has an ethylene polymerization activity similar to that of (L.sub.2 ZrCH.sub.3).sup.+ (C.sub.2 B.sub.9 H.sub.11).sup.-, whereas similar thorium complexes, e.g., (L.sub.2 ThCH.sub.3).sup.+ (Fe(C.sub.2 B.sub.9 H.sub.11)).sub.2.sup.-, are completely inactive for catalyzing ethylene polymerization due to strong Th-H-B interactions.\nIn general, both the [B(C.sub.6 H.sub.4 R).sub.4 ].sup.- type anions and the carborane anions exhibit a significant degree of coordination to the cations, and this suppresses the catalytic activity.\nThe advantage for the use of perfluorinated tetraphenyl borate as the counteranion is that it can made the catalysts more stable, and meanwhile largely overcome strong anioncation interactions. This makes the olefin polymerization activity of the homogeneous characterizable metallocene catalysts comparable to the IVB transition metal complex/MAO catalyst systems and the heterogeneous catalyst systems. The difference in activity between the cations with the fluorinated anions and cations with the nonfluorinated anions is dramatic. For example, (Me.sub.5 C.sub.5).sub.2 ThCH.sub.3.sup.+ B(C.sub.6 F.sub.5).sub.4.sup.- has an ethylene polymerization activity 3300 times higher than of (Me.sub.5 C.sub.5).sub.2 ThCH.sub.3.sup.+ B(C.sub.6 H.sub.5).sub.4.sup.-. An .sup.19 F NMR spectroscopic study shows that the four phenyl groups on the borane are equivalent at temperatures as low as -70.degree. C. in toluene-d.sub.8 while the X-ray structure shows a weak coordinative/interaction between the cation and m- and o-fluorine atoms from one of the four phenyl groups. In the case of zirconium, the difference in activity is around several hundred times.\nA wholly different and convenient way to generate the cationic species has also been reported. The neutral tris(pentafluorophenyl)borane can directly abstract a methide or hydride ligand from the neutral zirconocene complexes and thereby generate a cation and an anion. The NMR spectrum of (L.sub.2 ZrCH.sub.3).sup.+ (CH.sub.3 B(C.sub.6 H.sub.5).sub.3.sup.-, where L in this case is 1,3-dimethylcyclopentadienyl shows two signals for the methyl groups on the cyclopentadienyl rings in toluene-d.sub.8, indicating that the methyl groups are diastereotopic and the compound is a tight ion pair. The X-ray structure shows a weak linkage between the anion and cation through the abstracted methyl. This reaction also serves as a model for the generation of the cationic active centers in the heterogeneous system. The ethylene polymerization activity is 4.5.times.10.sup.6 gPE/mol.atm.h. (Nt=45s.sup.-1) at 25.degree. C.\nTo summarize, base-free group IVB or thorium cationic metallocene species are the active centers for olefin polymerization. The more open, both sterically and coordinatively, the more active the cation is. And the stability and noncoordinativity of anions are the key to keeping the cations stable and open."}
{"text": "Internal combustion engines are typically coupled to an emission control device known as a three-way catalytic converter (TWC) designed to reduce combustion by-products such as carbon monoxide (CO), hydrocarbon (HC) and oxides of nitrogen (NOx). Engines can operate at air-fuel mixture ratios lean of stoichiometry, thus improving fuel economy. For lean engine operation, an additional three-way catalyst commonly referred to as a Lean NOx Trap (LNT), is usually coupled downstream of an upstream three-way catalyst. The LNT stores exhaust components such as, for example, NOx and oxygen, when the engine is operating at a lean air-fuel ratio, and releases and reduces (purges) them when the engine is operating at a rich or stoichiometric air-fuel ratio. Over time, the ability of the LNT to store exhaust components can decrease due to factors such as sulfur deposits (SOx) from the fuel. Therefore, when the LNT efficiency is sufficiently reduced, a SOx purge has to be performed. Typically, the catalyst is heated and engine air-fuel ratio is changed to rich for SOx release and reduction. Since SOx purges result in fuel economy penalties, it is desirable not to purge unnecessarily. Thus, in order to maintain adherence to emission standards and obtain fuel economy benefits of a lean burning engine, it is desirable to monitor the efficiency of the LNT.\nOne method of determining the efficiency of the LNT is by correlating it to the oxidant storage capability, as described in U.S. Pat. No. 5,713,199. The amount of stored oxidant is calculated from a quantity of fuel required to purge it from the LNT wherein the purge is performed during high load operating conditions.\nThe inventors herein have recognized a disadvantage with this approach. Namely, when estimates of this method are performed under high load operating conditions, i.e., at high space velocity, the amount of time that the reductant used to purge stored oxidants is present in the LNT is reduced. Therefore, not all of the oxidants are purged during high-speed engine operation, resulting in an inaccurate estimate of the total LNT storage capacity. In other words, the reductant does not have sufficient time to react with stored oxidants when the LNT purge is performed under high load conditions. This will cause inaccurate capacity estimates and reductant wastage."}
{"text": "A computer network is a collection of computers and network hardware interconnected by communications channels that allow the computers and network hardware to create data connections across the communications channels and share data. The communications channels may be physically embodied in one or more wires, such as coaxial cables or fiber-optic cables, which connect the computers and network hardware, or the communications channels may be embodied using wireless communications protocols.\nData shared between computers may include documents, videos, and webpages, such as webpages for news websites or social networking websites. The data may allow users of the computers to collaborate on various projects or share information."}
{"text": "1. Field of the Invention\nThe present invention concerns a composite molding comprising a silicone gel molding and a peelable film, and a method for manufacturing the same.\n2. Description of the Related Art\nSilicone gel moldings are superior in terms of vibration-absorbing properties, electrical characteristics, heat resistance and water resistance, etc., and are therefore used in a broad range of applications such as shock-absorbing materials, protective materials, sealing materials and soundproofing materials, etc.\nHowever, since silicone gel moldings of this type have a strong surface tack, problems such as the adhesion of dirt and dust to the surfaces of such molding, the adhesion of such moldings to each other, and deformation resulting from the adhesion of such moldings to other materials, etc., have been encountered. Furthermore, since silicone gels themselves have a low mechanical strength, such gels are easily deformed and damaged.\nNumerous methods have been proposed in the past as methods for solving such problems. For example, methods have been proposed in which an organohydrogenpolysiloxane containing silicon-bonded hydrogen atoms is applied directly to the surface of a heat-cured silicone gel molding and diffused into the surface of said silicone gel molding, after which this organohydrogenpolysiloxane is heat-cured, thus forming a hard coating film on the surface of the silicone gel molding (see Japanese Patent Application Kokoku No. 1-25704 and Japanese Patent Application Kokoku No. 6-45222). However, although the surface tack problems of silicone gel moldings obtained using such methods is ameliorated to some extent, these methods do not completely solve the abovementioned problems of adhesion of dirt and dust, adhesion of moldings to each other and deformation of moldings. In the case of such methods, furthermore, two heating processes are required; accordingly, such methods are inferior in terms of productivity, and are not satisfactory as methods for the mass production of silicone gel moldings.\nOther methods of addressing the problems of dirt pick up, adhesion and deformation of silicone gel moldings are also known. In Japanese Patent Application Kokai No. 61-277414, a silicone molding in which the degree of cure varies in the direction of thickness is obtained by applying an organohydrogenpolysiloxane to the surface of a silicone gel composition which is capable of forming a silicone gel by being cured and then heating the resulting material before diffusion takes place. In Japanese Patent Application Kokai No. 5-69511 and Japanese Patent Application Kokai No 5-69512, a silicone gel molding which has a surface layer consisting of a silicone elastomer or silicone resin layer is obtained by causing a composition which forms a silicone elastomer or silicone resin layer after curing to contact the surface of a silicone gel composition, and then curing this assembly by heating. In Japanese Patent Application Kokai No. 6-88281, a silicone gel molding which has a high-hardness surface layer is obtained by pouring a silicone gel into a mold whose inside surface has been coated beforehand with an organohydrogenpolysiloxane, and then curing this assembly by heating. However, silicone gel moldings obtained using these methods suffer from the following problem: specifically, the components which form the surface layer are drawn into the silicone gel composition so that the curing of the surface layer becomes nonuniform. Furthermore, although the surface tack of silicone gel moldings obtained by the abovementioned methods is ameliorated to some extent, the abovementioned problems of adhesion of dirt and dust, adhesion of moldings to each other and deformation of moldings remain unsolved.\nIt is an object of the present invention to provide (a) a composite molding consisting of a silicone gel molding and a peelable film which does not suffer from the problems of adhesion of dirt and dust, adhesion of moldings to each other or deformation of moldings, and which is superior in terms of handling characteristics during storage, shipping or mounting on other materials, and (b) a method for manufacturing such a composite molding."}
{"text": "1. Technical Field\nThe present invention relates in general to data processing systems and, in particular, to a data processing system, method, and program for automatically testing software applications. Still more particularly, the present invention relates to a data processing system, method, and program for automatically generating a job within an automated test environment for testing software applications utilizing an event-driven work flow manager.\n2. Description of the Related Art\nPersonal computer systems are well known in the art. They have attained widespread use for providing computer power to many segments of today\"\"s modern society. Personal computers (PCs) may be defined as a desktop, floor standing, or portable microcomputer that includes a system unit having a central processing unit (CPU) and associated volatile and non-volatile memory, including random access memory (RAM) and basic input/output system read only memory (BIOS ROM), a system monitor, a keyboard, one or more flexible diskette drives, a CD-ROM drive, a fixed disk storage drive (also known as a xe2x80x9chard drivexe2x80x9d), a pointing device such as a mouse, and an optional network interface adapter. One of the distinguishing characteristics of these systems is the use of a motherboard or system planar to electrically connect these components together. Examples of such personal computer systems are IBM\"\"s PC 300 series, and Aptiva series.\nAn important part of software development is testing whether a particular software application functions as intended and without encountering errors. Typically, a large number of tests will be repeatedly executed on the software application. In order to perform these tests efficiently, a variety of software testing devices have been disclosed. These devices typically describe a testing method to be executed utilizing a single computer system which is also executing the software application being tested. The computer system will include the necessary test cases, any applications needed to execute the test cases, and all other necessary hardware or software components.\nA tester who desires to execute the tests must oversee and control the testing environment. The tester is required to define the test cases, identify the computer system to execute the test cases, indicate to the computer system when a particular portion of the test should start executing, when a portion of the test has completed and when to begin executing the next portion of the test. Although the computer system will automatically execute the tests, the tester must manually control the execution of the tests.\nA data processing system, method, and program including an automated software test environment are disclosed for generating a job to be executed within the automated software test environment to test a software application. A work flow manager is established for automatically managing the automated software test environment. The automated software test environment includes multiple computer systems coupled to a server computer system utilizing a network. The work flow manager is executed utilizing the server computer system. Execution of the automated test environment is controlled utilizing the work flow manager in response to a receipt of external events generated by the plurality of computer systems. A job description is created including an identification of one of the external events. A job described by the job description is executed utilizing the work flow manger in response to a receipt of the external event.\nAll objects, features, and advantages of the present invention will become apparent in the following detailed written description."}
{"text": "1. Field of the Invention\nThe invention in general relates to the fabrication of layered superlattice materials, and more particularly to a fabrication method using unreactive gas annealing and a low temperature pretreatment to reduce exposure of an integrated circuit to oxygen at elevated temperature.\n2. Statement of the Problem\nFerroelectric compounds possess favorable characteristics for use in nonvolatile integrated circuit memories. See Miller, U.S. Pat. No. 5,046,043. A ferroelectric device, such as a capacitor, is useful as a nonvolatile memory when it possesses desired electronic characteristics, such as high residual polarization, good coercive field, high fatigue resistance, and low leakage current. Layered superlattice material oxides have been studied for use in integrated circuits. U.S. Pat. No. 5,434,102, issued Jul. 18, 1995, to Watanabe et al., and U.S. Pat. No. 5,468,684, issued Nov. 21, 1995, to Yoshimori et al., describe processes for integrating these materials into practical integrated circuits. Layered superlattice materials exhibit characteristics in ferroelectric memories that are orders of magnitude superior to alternative types of ferroelectric materials, such as PZT and PLZT compounds.\nIntegrated circuit devices containing ferroelectric elements with layered superlattice materials are currently being manufactured. A typical ferroelectric memory cell in an integrated circuit contains a semiconductor substrate and a metaloxide semiconductor field-effect transistor (xe2x80x9cMOSFETxe2x80x9d) in electrical contact with a ferroelectric device, usually a ferroelectric capacitor. A ferroelectric memory capacitor typically contains a thin film of ferroelectric metal oxide located between a first, bottom electrode and a second, top electrode, the electrodes typically containing platinum. Layered superlattice materials comprise metal oxides. In conventional fabrication methods, crystallization of the metal oxides to produce desired electronic properties requires heat treatments in oxygen gas at elevated temperatures. The heating steps in the presence of oxygen are typically performed at a temperature in the range of from 800xc2x0 C. to 900xc2x0 C. for 30 minutes to two hours. As a result of the presence of reactive oxygen at elevated temperatures, numerous defects are generated in the single crystal structure of the semiconductor silicon substrate, leading to deterioration in the electronic characteristics of the MOSFET. Good ferroelectric properties have been achieved in the prior art using process heating temperatures at about 700xc2x0 C. to crystallize layered superlattice material. See U.S. Pat. No. 5,508,226, issued Apr. 16, 1996, to Ito et. al. Nevertheless, the annealing and other heating times in the low-temperature methods disclosed in the prior art are in the range of three to six hours, which may be economically unfeasible. More importantly, the long exposure time of several hours in oxygen, even at the somewhat reduced temperature ranges, results in oxygen damage to the semiconductor substrate and other elements of the CMOS circuit.\nAfter completion of the integrated circuit, the presence of oxides may still cause problems because oxygen from the thin film tends to diffuse through the various materials contained in the integrated circuit and combine with atoms in the substrate and in semiconductor layers forming oxides. The resulting oxides interfere with the function of the integrated circuit; for example, they may act as dielectrics in the semiconducting regions, thereby virtually forming capacitors. Diffusion of atoms from the underlying substrate and other circuit layers into the ferroelectric metal oxide is also a problem; for example, silicon from a silicon substrate and from polycrystalline silicon contact layers is known to diffuse into layered superlattice material and degrade its ferroelectric properties. For relatively low-density applications, the ferroelectric memory capacitor is placed on the side of the underlying CMOS circuit, and this may reduce somewhat the problem of undesirable diffusion of atoms between circuit elements. Nevertheless, as the market demand and the technological ability to manufacture high-density circuits increase, the distance between circuit elements decreases, and the problem of molecular and atomic diffusion between elements becomes more acute. To achieve high circuit density by reducing circuit area, the ferroelectric capacitor of a memory cell is placed virtually on top of the switch element, typically a field-effect transistor (hereinafter xe2x80x9cFETxe2x80x9d), and the switch and bottom electrode of the capacitor are electrically connected by a conductive plug. To inhibit undesired diffusion, a barrier layer is located under the ferroelectric oxide, between the capacitor\"\"s bottom electrode and the underlying layers. The maximum processing temperature allowable with current barrier technology is in the range of from 700xc2x0 C. to 750xc2x0 C. At temperatures above this range, the highest-temperature barrier materials begin to degrade and lose their diffusion-barrier properties. On the other hand, the minimum feasible manufacturing process temperatures of layered superlattice materials used in the prior art is about 800xc2x0 C., which is the temperature at which deposited layered superlattice materials are annealed to achieve good crystallization. Lower annealing temperatures require much longer time periods of exposure to oxygen, which can result in damage to the integrated circuit.\nFor the above reasons, therefore, it would be useful to have a method for fabricating layered superlattice materials in ferroelectric integrated circuits that minimizes the time of exposure to oxygen at elevated temperature, as well as reduces the maximum temperature used.\nThe embodiments of the present invention reduce the time of exposure of the integrated circuit to oxygen gas at elevated temperature, and reduce fabrication processing temperatures.\nIn an important embodiment of the invention, a portion of the time period during which an integrated circuit is heated or annealed at elevated temperature is conducted in an oxygen-free unreactive gas. The oxygen-free gas may be any relatively unreactive gas or mixture of unreactive gases, such as nitrogen and the noble gases, in particular, argon and helium. A useful result of embodiments of the invention is that when a fabrication method includes annealing in unreactive gas for a significant portion of the total annealing time at elevated temperature, then the ferroelectric polarizability of layered superlattice material is as high or higher than the polarizability of layered superlattice material annealed for the same total annealing time in oxygen only.\nThe invention provides a method of fabricating a thin film of layered superlattice material comprising: providing a substrate and a precursor containing metal moieties in effective amounts for spontaneously forming a layered superlattice material; applying the precursor to the substrate; annealing the solid thin film in an unreactive gas at a temperature in a range of from 600xc2x0 C. to 800xc2x0 C.; and annealing the solid thin film in an oxygen-containing gas at a temperature in a range of from 600xc2x0 C. to 800xc2x0 C. The invention contemplates that the annealing in an unreactive gas may either precede or come after the annealing in an oxygen-containing gas. The annealing in an unreactive gas may be conducted for a time period in the range of from 30 minutes to 100 hours. The annealing in an oxygen-containing gas is conducted for a time period in the range of from 30 minutes to two hours. In a preferred embodiment, the annealing in an oxygen-containing gas is conducted for a time period not exceeding 60 minutes. Typically, this step of annealing is conducted in substantially pure O2 gas.\nA further embodiment of the invention includes heating the substrate after applying the precursor, forming a solid film from the precursor. In an embodiment, the heating is conducted at a temperature not exceeding 600xc2x0 C. The heating typically comprises a step of drying the precursor coating on the substrate at a temperature not exceeding 300xc2x0 C. Typically, drying is accomplished by baking in an oxygencontaining gas, preferably in O2 gas.\nIn another important embodiment of the invention, the heating comprises a step of pretreating the substrate after the precursor is applied and before annealing. Typically, pretreating is conducted in an oxygen-containing gas. Pretreating preferably comprises rapid thermal processing (xe2x80x9cRTPxe2x80x9d) of the substrate having the precursor coating. In a typical embodiment, the rapid thermal processing is conducted at a temperature in the range of from 300xc2x0 C. to 600xc2x0 C. for a time period in the range of from 1 minute to 15 minutes. Pretreating also may comprise a hot plate baking of the substrate. Typically, hot plate baking is conducted in air at a temperature in the range of from 300xc2x0 C. to 600xc2x0 C. for a time period in the range of from 1 minute to 15 minutes. Pretreating may also comprise a furnace pre-annealing of the substrate. Typically, furnace pre-anneal is conducted at a temperature in the range of from 300xc2x0 C. to 600xc2x0 C. for a time period in the range of from 1 minute to 15 minutes.\nIn another embodiment of the invention, the substrate comprises a first electrode, and the method includes steps of forming a second electrode on the solid thin film, after the step of annealing, to form a capacitor, and subsequently performing a step of post-annealing. In a preferred embodiment, the first electrode and the second electrode contain platinum. The step of post-annealing is conducted at a temperature not exceeding 800xc2x0 C., preferably for a time period in the range of from 30 minutes to two hours. In one embodiment of the invention, the post-annealing is conducted in an oxygen-containing gas, typically in O2 gas. In a preferred embodiment of the invention, the post-annealing is conducted in oxygen-free unreactive gas, typically N2 gas.\nIn a preferred embodiment of the invention, an electrically conductive barrier layer is formed on the substrate prior to applying the precursor coating.\nThe thin film of layered superlattice material typically has a thickness not exceeding 500 nanometers (xe2x80x9cnmxe2x80x9d). In one embodiment, the layered superlattice material comprises strontium, bismuth and tantalum. In another embodiment, the precursor includes metal moieties for forming a layered superlattice material thin film comprising strontium, bismuth, tantalum and niobium.\nIn accordance with the invention, the precursor comprises a solution of metal organic precursor compounds containing metal atoms contained in the desired layered superlattice material. Preferably, the metal organic compounds are metal 2-ethylhexanoates.\nNumerous other features, objects and advantages of the invention will become apparent from the following description when read in conjunction with the accompanying drawings."}
{"text": "In unsafe environments and particularly explosive environments, such as a mine environment or an oil well environment, it is necessary that the amount of energy dissipated into the surrounding atmosphere from electronic components disposed in the unsafe environment needs to remain below certain levels (e.g. 300 micro Joules) to avoid igniting a mixture of methane and air that would result in an explosion. In order to ensure that any electronic components disposed in the unsafe environment are not a risk for causing an explosion, safe area electronic systems used to drive the electronic component are required to output a signal that is guaranteed to stay below certain voltage and current levels."}
{"text": "This invention relates to a display material which is capable of reversibly changing its coloring condition by application of heat and an electric field, and more particularly to a novel display material which is useful as electrochromic display element material and thermosensitive recording material.\nFor example, in the field of thermosensitive recording paper, there has been used a display material which has a leuco-dye dispersed in the form of crystallite in a binder together with a developer or activator so that, when thermally fused, the pigment and the developer are contacted with each other on a molecular level to develop a color. More specifically, the conventional thermosensitive recording paper which is in wide use has crystallites of a leuco dye and phenolic compound as the developer dispersedly coated together with a hydrophilic binder on substrate paper (a support body) such as synthetic paper or the like.\nHowever, in consideration of the manufacturing process, the conventional thermosensitive recording paper of such a construction unavoidably suffers from the problem of fogging (a slight degree of coloration before printing). Namely, it is difficult to avoid a certain degree of reaction between the crystallites of the leuco dye and phenolic compound as the developer during the manufacturing process in the stages of kneading them into a hydrophylic binder, coating and drying, as a result coloring the heat sensitive paper as a whole in grey to a certain degree before serving for printing. Thus, at the present moment it is difficult to obtain thermosensitive recording paper which is satisfactory in whiteness. In addition, the above-mentioned conventional thermosensitive recording paper is susceptible to colorations by inadvertent heating, for example, colorations by application of heat during storage or after printing, making the printed information illegible or inviting considerable degradations in the quality of picture images.\nNevertheless, due to the steady increase in the amount of documents which are handled in the office work, there have been strong demands for thermosensitive recording paper which can be used repeatedly, from a viewpoint of saving space and resources. Besides, polychromic coloration of thermosensitive paper is considered for improving visibility of documents or the like.\nHowever, the demands cannot be fulfilled by the conventional thermosensitive recording paper, including the sublimation transfer type and the thermal fusion type which are also unsatisfactory in case of high-speed full-color image production by a thermal head and are almost impossible to use repeatedly for reproduction.\nOn the other hand, in the field of the display, researches have been made for development of a display material which is suitable for the so-called flat display of a flat plate-like shape, including development of electrochromic display elements.\nIn this connection, the conventional electrochromic display elements are mostly designed for monochromic coloration and decoloration, and scarcely capable of coloration and decoloration in three primary colors such as RGB (the three primary colors in the additive color method) and YMC (the three primary colors in the subtractive color method) suitable for the full-color display.\nSince, as mentioned hereinbefore, the conventional display materials used in the fields of the thermosensitive recording paper and display invariably lack stability and are difficult to apply to repeated use and polychromic coloration and insufficient in color vividness, the development of a novel display material which can eliminate these problems is greatly expected."}
{"text": "With the popularity and wide access of broadband networks, software publishers have made applications available for use on central servers. Indeed, a commercial market segment known as Application Service Providers (ASPs) has recently emerged. Broadband network access facilitates the use (including downloading and execution) of such server-based software applications by providing high speed data connections for distribution of software applications from a server terminal (or from a remote location) to a local client terminal (for example, a personal computer terminal). However, many users consider it troublesome and time consuming to install applications distributed in this manner onto their client terminals because it typically requires accessing the server terminal via the data connection, downloading the desired application, and thereafter installing the downloaded application onto the client terminal for execution.\nIn view of the foregoing, it is desirable to have a system and method that provides unobtrusive execution of externally delivered applications, without the traditional installation procedure of the application before execution. It is further desired to have secure access to the applications to prevent unauthorized users from executing them."}
{"text": "In the perfumery as well as in the flavour industry a lot of efforts are devoted today to the replacement of expensive naturally occurring raw materials or the reproduction of original organoleptic effects by making use of new chemicals. Among the numerous synthetic odorants prepared in a recent past, bi-, tri- or polycyclic derivatives of decalin can be cited as examples, as well as bicyclo[2.2.1]heptane, bicyclo[2.2.2]octane, cedrene or caryophyllene derivatives.\nFor instance, the tricyclic compound of formula ##STR2## is appreciated in the art for its original floral and woody odour--see e.g. Swiss Pat. No. 547 850--, as well as the compound of formula ##STR3## which is used in perfumery for developing or enhancing woody and spicy odour notes--see e.g. Swiss Pat. No. 557 870--. As further example of tricyclic fragrant compounds, one can also cite the compound of formula ##STR4## characterized by its original odour note, at the same time sweet, fresh, spicy and herbaceous--see e.g. DE-OS No. 28 26 302.\nNotwithstanding the relative abundance of tricyclic fragrant compounds known in the art, there still remains a constant need of enrichment of the palette of the perfumer, or of the flavourist, by the finding of new active substances."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for controlling the operation of a hard disk drive device based upon a detected value of surrounding atmospheric pressure and a hard disk drive device capable of controlling its operation based upon the detected value of the atmospheric pressure.\n2. Description of Related Art\nAt least one magnetic recording disk is contained in a hard disk drive device, and the data is stored on both surfaces of the magnetic recording disk. One read/write head is provided on each recording surface. In one example, four read/write heads are used with two recording disks; each read/write head is mounted on a front end of a head support arm, and rear ends of four head support arms are so connected to each other that all the read/write heads address the same radial cylinder positions on each of the four recording surfaces of the two recording disks. In this manner, all the read/write heads address the same radial positions of the four recording surfaces which define a cylinder, and as such the radial cylinder positions are called cylinder position.\nA prior hard disk drive device used a contact start/stop scheme in which a head/slider assembly is landed on a rest region or a non-recording region located inside the inner most recording cylinder of a magnetic recording disk during a power off period, and takes off from the surface of the rest region when the magnetic recording disk is rotated. To realize a reliable take off, the surface of the disk must be roughened. The reasons for forming the roughness on the recording surface is that if the surface is made smooth, the head/slider assembly sticks to the recording surface due to interatomic forces, thereby preventing take off.\nAs the diameter of the magnetic recording disk becomes smaller, many efforts have been made to increase the recording density and to improve S/N ratio.\nTo increase the recording density and to improve S/N ratio, a flying height of the head/slider assembly, i.e., a space between the head/slider assembly and the recording surface of the disk has been decreased. In the magnetic recording disk, it was difficult to decrease the height of the head/slider assembly over the roughened recording surface, since the head/slider assembly tended to contact to top portion of the roughened surface.\nA hard disk drive device has been developed which uses a load/unload scheme in which a ramp element 2 is mounted at the peripheral of the magnetic recording disk 1, as shown in the FIGS. 1(A) and 1(B). The surface of the ramp element 2 is ramped, and a front end 3 of a head support arm 4 rides on the ramped surface of the part 2, and moves in the direction of an arrow 6 to the rest position of the ramp element 2, as shown in FIG. 1(B), when the head support arm 4 is rotated around a pivot point 5 in the clockwise direction in FIG. 1(A). When the read/write operation is started, the magnetic recording disk 1 is rotated by a motor, not shown, and the head support arm 4 is moved from the rest position in a direction of an arrow 8, and is moved to a flying position above the desired cylinder of the magnetic recording disk 1 to read the data from the cylinder or write the data into the cylinder.\nThe rotational speed, i.e., the revolution per minute (RPM), of the magnetic recording disk 1 is so designed to generate an air bearing with an appropriate pressure which causes the head/slider assembly 7 of the head support arm 4 to fly above the surface of the magnetic recording disk 1.\nIn this manner, the head/slider assembly 7 does not land on the surface of the disk 1, so that the surface of the disk 1 can be polished to remove the roughness of the surface of the disk used in the contact start/stop scheme, and the flying height of the head/slider arm 4 in the read/write mode and an error recovery mode can be reduced in comparison with the flying height of the contact start/stop scheme.\nDescribing the flying height of the head/slider assembly 7 in the read/write mode and the error recovery mode with reference to FIG. 2, in the read/write mode, the flying height of the head/slider assembly 7 is maintained at the height H1 by maintaining the first rotation speed. If the hard disk drive device senses a read error or a write error, the hard disk drive device enters the error recovery mode by reducing the rotational speed of the disk 1 to a second rotational speed which is lower than the first rotational speed to maintain the height H2. This height H2 is so selected to bring the lower surface of the head/slider assembly 7 close to the surface of the disk 1 to wipe off any dusts or residual material on the surface of the magnetic recording disk 1.\nThe problem found by the inventors of the present invention is that the head/slider assembly 7 tends to stick to the surface of the magnetic recording disk 1 when the hard disk drive device enters the error recovery mode while the hard disk drive device is used in an airplane flying at a high altitude."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may or may not constitute prior art.\nLithotripsy is a common method for fragmenting stones, or calculi, in the urinary tract, kidneys, and/or bladder. Lithotripsy, however, can leave stone fragments, which act as nucleation sites for future stone formation, for example, in the lower pole of a kidney. These fragments remain in the patient after the procedure mostly because of the difficulty in capturing very small fragments of dust employing conventional removal devices such as, for example, stone baskets.\nTypical corrective actions may include one or more of the following: monitoring for future stone formation, performing additional ureteroscopic procedures, or percutaneous nephrolithotomy. Among the literature that can pertain to this technology include the following patent documents and published patent applications: EP 2355717, US 20090136594, JP 2006314811, U.S. Pat. Nos. 7,883,516, and 8,753,351, the entire contents of which are incorporated herein by reference for all purposes.\nAccordingly, there is a need for a device that can be employed with lithotripsy that removes small fragments. It would be desirable to have mechanisms which attract, retain, and contain stone fragments to reliably remove small stone fragments and debris."}
{"text": "Spinal fusion is a common technique used to treat chronic back pain caused by degenerated or herniated disk. The technique involves the removal of a disc between two vertebrae and replacing it with an intervertebral spacer. The intervertebral spacer maintains spacing between the two vertebrae and preferably results in fusion through the spacer. The intervertebral spacers may be constructed of autogenic bone tissue taken from a patient's own bone. Allogenic spacers are constructed of bone harvested from donors. Artificial spacers are currently the most common spacer type and may be constructed of metallic material such as titanium or stainless steel or polymers such as polyetheretherketone (PEEK).\nPEEK has recently become popular due to its biocompatibility and naturally radiotranslucent characteristics, resulting in limited interference with x-ray and CT imaging. However, while PEEK is biocompatible, bone treats it as a foreign body during the remodeling process and isolates it with a fibrous tissue capsule. This fibrous tissue prevents direct bony apposition and adhesion to the implant. Other materials, such as titanium, allow for direct bony apposition and ongrowth, but they are typically not radiotranslucent and it becomes difficult to assess the fusion formation.\nOther areas where PEEK is used as an orthopedic biomaterial experience similar fibrous encapsulation. Such indications include custom machined bodies that are used to fill defects in the skull and cranium. With PEEK, MRI and CT imaging is generally easier as compared to titanium, but the implant is never fully incorporated into the bone and soft tissue does not adhere to the implant.\nCeramic materials such as calcium phosphates, β-TCP, hydroxyapatites and the like allow for direct bony apposition much like titanium. However, they are typically limited in their strength and toughness. Therefore, it is desirable to construct a material that combines more of the desired properties from other individual materials described above, such as toughness and strength, less interference with MRI, X-ray or CT imaging, tissue adhesion, etc."}
{"text": "DC to DC converters are electronic circuits that convert a source of direct current from one voltage to another. Switching DC to DC converters are commonly used because of their higher efficiency, which is especially important when the electronic device connected to the DC to DC converter is a battery-powered portable device.\nFIG. 1 illustrates an example DC-DC converter 100 based on a step-down switching circuit particularly suitable for meeting many of the general goals for supplying power to advanced digital circuits. Circuit 100 in FIG. 1 (simplified to facilitate illustration and discussion) is called a DE-DRC (Differentially Enhanced-Duty Ripple Control) circuit. The DE-DRC design has been previously published by the present inventor and others. See, for example, J. Fan, X. Li, S. Lim, and A. Huang, “Design and Characterization of Differentially Enhanced Duty Ripple Control (DE-DRC) for Step-Down Converter,” IEEE Trans. Power Electron., vol. 24, no. 12, pp 2714-2725, December 2009. In FIG. 1, a power source 102 provides direct current at an input voltage VIN.\nThe circuit 100 provides direct current to a load (RLOAD) at an output voltage VOUT. Two electronic switches (SW1, SW2) are controlled by a switch control circuit 110 and driver 112. At most only one switch is closed at any one time. A comparator 108 controls the switch control circuit 110. There are two feedback paths. In a first feedback path, the voltage VSW on the switched side of the inductor LO is coupled to the comparator 108 through a low pass filter (RR, CR). In a second feedback path, two differential difference amplifiers (104, 106) generate a differential pair of feedback signals, VP and VN. These two differential signals are coupled to comparator 108, where VN is directly coupled, and VP is coupled through a high pass filter (CR, RR).\nVP and VN are as follows:VP=KP(HVOUT−VREF)+VOUT VN=−KN(HVOUT−VREF)+VOUT Where H=RS2/(RS1+RS2), KP=gain of positive differential circuit 104, and KN=gain of negative differential circuit 106, and VREF is a constant reference voltage.\nThe first feedback path coupling VSW to the comparator 108 is fast because there are only passive components between VSW and the comparator. The feedback path coupling VOUT to the comparator 108 is slower, because VOUT is proportional to the integral of current in COUT, and there are active amplifiers between VOUT and the comparator. The second loop can adjust the control bandwidth. Both loops have a big influence on the transient response. The two feedback loops combined provide a stable system with high bandwidth control.\nTo simplify the discussion, assume H=1 and VIN is constant. First, consider the steady state (constant load). In the steady state, with H=1, VCONTROL (the positive input of the comparator 108) is approximately VOUT. VSW is a square wave, having an average value of approximately VOUT, and low-pass filtered VSW contributes a sawtooth waveform having an average value of approximately the average value of VSW to VRIPPLE (the negative input of the comparator 108). When VRIPPLE drops below VCONTROL, comparator 108 causes switch control circuit 110 to close SW1 for a constant on-time. VSW is driven higher while SW1 is closed, and VOUT is controlled to be approximately equal to the average value of VSW. If the load current increases, the average value of VSW increases rapidly (because of increasing current through LO), VP has a transient decrease (VOUT decreases due to an increased current draw from COUT), VN has a transient increase, and the duty cycle of SW1 is increased (switching frequency increases). If the load current decreases, the average value of VSW rapidly decreases, VP has a transient increase, VN has a transient decrease, and the duty cycle of SW1 is decreased. The magnitude of VRIPPLE at the negative input of the comparator 108 is relatively large compared to the ripple voltage on VOUT, which provides good noise immunity in the feedback signal. In addition, the circuit provides a fast response to load transients over a wide input and output range.\nFIG. 2 shows an improved controller disclosed in co-pending commonly-assigned application Ser. No. 13/171,283 having common inventorship with the present application (see FIG. 3). This application is incorporated herein by reference in its entirety and for all purposes.\nIn FIG. 2, there are three feedback paths. In a first feedback path, the voltage VSW on the switched side of the inductor LO is AC coupled to the comparator 208. In a second feedback path, two differential difference amplifiers (204, 206) generate a differential pair of feedback signals, VP and VN. These two differential signals are coupled to comparator 208, where VN is directly coupled, and VP is coupled through a high pass filter network (CR1, RR)(CR2, RDC). In a third feedback loop, the reference voltage is adjusted by a reference adjust circuit 214.\nVP and VN are as follows:VP=KP(HVOUT−VREF adjusted)+VBIAS VN=−KN(HVOUT−VREF adjusted)+VBIAS Where H=RS2/(RS1+RS2), KP=gain of positive differential circuit 204, and KN=gain of negative differential circuit 206, VREF adjusted is a variable reference voltage, and VBIAS is a fixed reference voltage.\nThe first feedback path coupling VSW to the comparator 208 is fast, because VRIPPLE is proportional to AC current in inductor LO, and there are only passive components between VSW and the comparator. The feedback path coupling VOUT to the comparator 208 is slower, because VOUT change is proportional to the integral of AC current in COUT, and there are active amplifiers between VOUT and the comparator. The feedback path adjusting VREF is intentionally very slow to avoid interference with the first two feedback paths. The third feedback path provides better output voltage accuracy. The three feedback paths combined provide a stable system with high bandwidth control and accurate output voltage.\nThese two control circuits are well-suited for controlling a single phase DC to DC converter power supply. However, there is a growing need for high amperage power supplies. For example, if one needs a 180 amp power supply and has an existing design for a 30 amp power supply, this need may be met by combining 6 of the 30 amp power supplies in parallel. This saves the cost of redesigning the power supply for higher amperage but it creates a problem which will be illustrated below in connection with FIGS. 3 and 4.\nThis problem is most easily explained utilizing a multiphase power supply comprising two DC to DC converters connected in parallel (phases 1 and 2, respectively) and having a single output terminal connected to a load. FIG. 3 illustrates the waveforms, generally indicated as 300, for the case where this two-phase system is operating at a 20% duty cycle. The current signal sensed for phase 1 is shown in FIG. 3A as 302. The current signal sensed for phase 2 is shown as 304 in FIG. 3B. The current sense signals 302, 304 are combined in FIG. 3C to yield ripple signal 306. In FIG. 3D, the ripple signal 306 is compared against a comp signal 308. When the ripple signal falls below the comp signal 308 (which is the output voltage of an error amplifier which compares a feedback voltage with a reference voltage), an ON pulse trigger signal such as 310, 312, 314, 316, 318 and 320 illustrated in FIG. 3F, is generated. The first ON pulse trigger signal 310 triggers the generation of a PWM signal 322 for phase 1, as shown in FIG. 3G. The second pulse trigger signal 312 triggers the generation of a PWM signal 324 for phase 2, as shown in FIG. 3H. The third ON pulse trigger signal 314 triggers the generation of PWM signal 326 for phase 1 (FIG. 3G). The fourth ON pulse trigger signal 316 triggers the generation of PWM signal 328 for phase 2 (FIG. 3H), and so on.\nFIG. 4 illustrates the waveforms of the same circuit as FIG. 3, generally indicated as 400, except that the duty cycle has been increased to 50%. The current signal sensed for phase 1 is shown in FIG. 4A as 402. The current signal sensed for phase 2 is shown as 404 in FIG. 4B. The current sense signals 402, 404 are combined in FIG. 4C to yield ripple signal 406. The problem is that if both phases have equal currents and are exactly 180° out of phase, the current signals 402, 404 will cancel out and no ripple signal would be available to control the converters. Therefore, when the signal 406 is compared to the signal comp, no ON pulse trigger signals are generated as shown at 408 in FIG. 4E. Therefore, no PWM signals are generated for both phases 1 or 2 as shown in FIGS. 4F and 4G. Thus, the multiphase power converter is out of control. Even if the currents to each of the DC to DC converters were not the same and/or were not exactly out of phase, the ripple signal that results would be very small so that noise signals in the millivolt range could cause erratic operation of the multiphase power converter\nAlthough the problem is most easily illustrated with a 2 phase multiphase power converter, the problem is common to any multiphase power converter system using combined current information for control. For a 3 phase system, the maximum duty cycle before this problem appears is 33.3%. For a 4 phase converter the maximum duty cycle before this problem appears is 25%. For the 6 phase converter discussed above, the maximum duty cycle before this problem appears is an unacceptable 16.6%. In general, the maximum duty cycle allowable before this problem appears is 1/n*100%, where n is the number of phases in the multiphase power converter system.\nTherefore, there is a need for a controller for a multiphase power converter in which the duty cycle is not so limited."}
{"text": "The current invention relates to an electronic component, in particular a chip element.\nIt is known to produce sensors that function in a contactless manner, e.g. for the detection of measurement quantities such as speed, angle, or path, by means of magnetic thin film technology. Worth particular mention in this regard are AMR thin film technology using the properties of anisotropic magnetoresistive sensor materials. The use of so-called GMR materials (giant magnetoresistance) has turned out to be particularly advantageous.\nBy using virtually identical technology, a generically related functional layer material is used to produce magnetic memory elements, so-called MRAMs (magnetic random access memory). In comparison to conventional memory elements, these MRAMs have the advantage that on the one hand, they can be very rapidly read and rewritten and on the other hand, they are in a position to retain their memory content without permanent power supply.\nIn the sensors described above, it turns out to be disadvantageous that the disposition of individual sensor elements on a substrate, which disposition is necessary for their function, takes up a lot of space. For example, in known AMR or GMR angle sensors, the functional layer is comprised of long, meandering strip conductors which must be laid exactly in relation to one another geometrically in order to receive a suitable measurement signal. This produces large intermediate spaces which makes the sensor as a whole large.\nTherefore, the object of the invention is an improved usage of the space occupied by a magnetoresistively functioning sensor.\nThe objects of the present invention are attained by an electronic component, particularly a chip element comprising first magnetoresistive elements disposed on a substrate and having a sensor function, and at least one second magnetoresistive elements having a memory function, said first magnetoresistive elements consisting of elongated meandering strip conductors which are disposed in relation to each other exactly geometrically to obtain signals to be measured, wherein intermediary spaces are formed between said first magnetoresistive elements having said sensor function on said substrate, and said second magnetoresistive elements having said memory function are arranged in said intermediary spaces between said first magnetoresistive elements having said sensor function.\nAccording to the invention, it is now possible to make sensible use of the large amount of space required in magnetoresistively functioning sensors. According to the invention, free areas between sensor components can be used for memory elements. In comparison to a separate production of the sensor and memory, different components such as housings can be saved, which can reduce costs.\nSuitably, the at least one first magnetoresistive element is an AMR sensor or element or a GMR sensor or element. As a rule, elements of this kind are comprised of long, meandering strip conductors which are disposed exactly in relation to one another geometrically. For example, in order to produce sine and cosine signals of an angle to be measured (by means of an external magnetic field associated with the angle to be measured) two AMR sensor elements are disposed at an angle of 45xc2x0 in relation to each other, wherein the angle associated with the sine and cosine signals can be determined by using the arc tan function. With a device of this kind, intermediary spaces are necessarily produced between the sensor elements on the chip surface or the substrate, which can be advantageously used according to the invention to form memory elements. Since the memory elements have significantly smaller dimensions than the sensor elements, the available area can be optimally utilized.\nAccording to a preferred embodiment of the invention, the at least one second magnetoresistive element is an MRAM cell. Cells of this kind have particularly small dimensions and can be simply disposed between sensor elements on the substrate. MRAM cells of this kind have the same layer structure as certain AMR or GMR sensor elements so that it is particularly easy to equip a substrate in this instance.\nIn a suitable fashion, the component according to the invention has an integrated set of control and evaluation electronics for the sensor function and/or the memory function. This produces a particularly compact component.\nAccording to an advantageous embodiment of the component according to the invention, the at least one second magnetoresistive element can store data for executing and/or controlling and/or adjusting the sensor function of the at least one first magnetoresistive element. For example in angle sensors, the memory area can be significantly increased by virtue of the fact that the number of revolutions occurring or measured can be stored in the integrated memory. A particular advantage is based on the fact that the power supply of this kind of sensor and memory elements can be interrupted without signal loss. The memory content and a current sensor signal are immediately available again when operation is resumed. In magnetoresistively functioning path sensors in which a magnetoresistive element scans a periodic magnetic field pattern, the number of periods detected can be stored in a memory and can be supplied to a set of evaluation electronics in a relatively rapid fashion.\nFurthermore, offsets that occur, for example when there are a number of magnetoresistively functioning sensors, are corrected by means of an expensive thin film compensation. According to the invention, offset data or compensation data can be stored in the integrated memory elements. Conventional measures or electronic compensation processes such as focal length calibration, zener zapping, thyristor zapping, etc. can be easily replaced. Moreover, arbitrary characteristic signal data can be stored in the memory elements. Through the use of MRAM memory cells, data of this kind can be quickly accessed and can be updated, for example, in order to adapt the functional parameters (offset, sensitivity, operating point) or in order to execute an adaptive signal correction or a calibration. The stored data can also be used within the framework of a self-diagnosis of the sensor."}
{"text": "In media handling assemblies, particularly in printing systems, accurate and reliable registration of the substrate media as it is transferred in a process direction is desirable. In particular, accurate registration of the substrate media, such as a sheet of paper, as it is delivered at a target time to an image transfer zone will improve the overall printing process. The substrate media is generally conveyed within the system in a process direction. However, often the position and/or timing of the substrate media can deviate from that which is intended or desired. The sheet might be ahead or behind in its process direction position, or the sheet can shift in a cross-process direction (lateral to the process direction) or even acquire an angular orientation (referred herein as “skew”) such that the opposed linear edges are no longer parallel to the process direction. Thus, there are three degrees of freedom in which the substrate media can move, which need to be controlled in order to achieve accurate delivery thereof. A slight skew, lateral misalignment or error in the arrival time of the substrate media through a critical processing phase can lead to errors, such as image and/or color registration errors relating to arrival at an image transfer zone. Also, as the substrate media is transferred between sections of the media handling assembly, the amount of registration error can increase or accumulate. A substantial skew and/or registration error can cause pushing, pulling or shearing forces to be generated, which can wrinkle, buckle or even tear the sheet.\nContemporary systems transport a sheet and deliver it at a target time to a “datum,” based on positional measurements from the sheet. That datum, also referred to herein as a delivery registration datum, can be a particular point in a transfer zone, a hand-off point to a downstream nip assembly or any other target location within the media handling assembly. Typically, the time and orientation of the sheet arriving in a sheet registration system is measured by sensors located near the input of the registration system. A controller then computes a sheet velocity command profile designed to deliver the sheet at a target time that delivery registration datum. A sheet velocity actuator commanded by the controller then executes a command profile in order to timely and accurately deliver the sheet. Examples of typical sheet registration and deskewing systems are disclosed in U.S. Pat. Nos. 5,094,442, 6,533,268, 6,575,458 and 7,422,211, commonly assigned to the assignee of record herein, namely Xerox Corporation, the disclosures of which are each incorporated herein by reference. While these systems particularly relate to printing systems, similar paper handling techniques apply to other media handling assemblies.\nSuch contemporary systems attempt to achieve position registration of sheets by separately varying the speeds of laterally spaced apart drive rollers in registration nip assemblies to correct for skew mispositioning of the sheet, which is also referred to as differentially driven drive or nip assemblies, such as that disclosed in U.S. Pat. No. 7,422,211. By imparting specific differential drive velocity profiles to the two drive nips over a small period of time, skew, process direction and/or lateral position of the sheet can also be corrected. Separate drive motors and/or belt assemblies are often included in differential drive systems, for imparting an angular velocity to the driven wheels. While each motor may be connected directly to the driven wheels, belts (also referred to as timing belts) are often employed. Also, the motors may be stepper motors or DC servo motors with encoder feedback from an encoder mounted on the motor shaft, a driven wheel shaft or the idler shaft. Such registration nip assemblies also generally includes sheet sensors, which are used to detect the arrival of a sheet, its lateral position, skew and other characteristics. Temporarily driving the laterally spaced nips at slightly different rotational speeds will produce a slight difference in the total rotation or relative pitch position of each drive roll while the sheet is held in the two nips. In this way, one side of the sheet moves ahead of the other to induce skew (small partial rotation) in the sheet, in order to eliminate and/or correct for detected skew or positional errors in the lateral or process directions.\nAlternatively, contemporary systems include a translating carriage on which the registration nip assemblies are mounted, such as that disclosed in U.S. Pat. No. 5,094,442. As shown in FIG. 6, a nip assembly 2 includes a driven wheel 6 (also referred to as a drive roll) and an idler wheel 8, (also referred to as an idler roll) which together engage opposed sides of the sheet S and conveying it within the printing system in a process direction P. The system includes two laterally spaced apart nip assemblies 2 that are together mounted on a carriage 40. The carriage 40 is able to translate laterally with the use of a separate motor 42 and screw drive shaft 44, as well as a carriage guide collars 46 slideable along a carriage guide shaft 48. The motor 42 turns the screw drive shaft 44, which then translates the carriage 40 laterally, along with the nip assembly 2. In this way, as the carriage 40 with the nip assembly 2 translates laterally, so does the sheet S.\nFurther sheet registration systems are disclosed in U.S. Pat. Nos. 5,697,608 and 6,866,260, commonly assigned to the assignee of record herein, namely Xerox Corporation, the disclosures of which are each incorporated herein by reference. Such systems use a pair of sheet edge sensors, located on one side of the sheet path, to measure the position of a sheet upon arrival in the sheet registration nip assembly. One of the two edge sensors is generally located laterally adjacent or just upstream of the registration nip assembly, with the other edge sensor disposed further upstream. In this way, when the sheet arrives at the registration nip assembly, the differential measurements from the two edge sensors can be used to calculate lateral position and skew of the sheet. This information is then fed to a controller, which in turn signals the registration nip assembly in order properly register the sheet position laterally and in skew. Typically, the controller calculates the correction of the sheet lateral and skew position to be completed prior to each sheet's arrival at the downstream delivery registration datum. That earlier point for completion of the registration correction is a virtual registration datum that lies somewhere between the registration nip assembly and the delivery registration datum. However, often a sheet can arrive at the virtual registration datum with its registration not fully corrected. Also, further registration errors can occur as the sheet travels from the virtual registration datum to the delivery registration datum.\nAccordingly, it would be desirable to provide a method and apparatus capable of more accurately registering a sheet in a media handling assembly, which overcomes the shortcoming of the prior art."}
{"text": "It is generally accepted that DNA is the crucial target in the cytotoxic effects of ionising radiation. There is considerable evidence to support the view that DNA double-stranded (ds) breaks are particularly important. The DNA damage results from both direct ionisation in the DNA molecule (direct effect) and by indirect effects mediated by the radiolysis products of water. Carbon-centred radicals on the deoxyribose moiety of DNA are thought to be important precursors of strand breaks. Ionising radiation also induces damage in DNA bases. If the level of cellular DNA damage is sufficient, the consequence of irradiation is cell kill, and thus ionising radiation is used as a mode of cancer therapy. For irradiated normal tissues, the cell killing can result in temporary or permanent impairment of tissue and organ function. The extent of such effects is dependant upon the radiation dose, and if sufficient, can be lethal to the organism. For humans and other animals, hematopoiesis is the most radiosensitive organ/function, followed by the gastrointestinal mucosa. Even if radiation induced DNA damage is sublethal, mutagenic lesions can have serious long term consequences, including carcinogenesis.\nThe medical strategies or countermeasures aimed at reducing the extent of radiation-induced effects are broadly described as radioprotectors (which to be effective, generally need to be administered prior to radiation exposure), mitigants/mitigators (which can be effective if administered after irradiation, but before the appearance of symptoms), and treatments which are generally administered after the appearance of symptoms. A sub-class of the prophylactic radioprotectors are drugs that reduce the extent of the initial radiation-induced DNA damage, and it is this sub-class that is the major focus of the present invention.\nThe commercial potential of radioprotectors resides primarily in two distinct arenas. One of these relates to the need to protect normal tissues in cancer radiotherapy patients, and the other concerns the need to assuage the consequences of unplanned irradiation associated with civil scenarios, such as radiation accidents and radiation terrorism, as well as irradiation in military contexts. This second scenario would also include planned exposure to ionising radiation in medical diagnostic procedures, in which administration of radioprotectors could ameliorate the health risks associated with low or modest radiation doses, without interfering with diagnostic imaging processes.\nThe treatment of tumours with ionising radiation (hereinafter referred to as “cancer radiotherapy”) is used extensively in cancer therapy. The goal of such treatment is the destruction of tumour cells and inhibition of tumour cell growth presumably through DNA damage, while minimising damage to non-tumour cells and tissues. The potential for damage to non-tumour cells in the vicinity of the tumour limits the radiation dose that can be administered, which in turn often limits the effectiveness of radiotherapy against certain tumours. This is especially the case in relation to brain tumours and tumours in the abdominal cavity.\nCancer radiotherapy is a very significant public health activity. Given the incidence of cancer in the population and the international assessment that more than 50% of cancer patients benefit from inclusion of radiotherapy in their treatment, more than 10% of the population are likely to experience cancer radiotherapy in their lifetime.\nThe dominant consideration in prescribing radiation doses for cancer radiotherapy is the assessment of tolerance of the most radiosensitive normal tissues/organs in the treatment field. This assessment, together with the expected radiation dose required to eradicate a tumour, often determines whether the treatment strategy is aimed at cure or palliation. In many cases, the maximum tolerable doses are insufficient to eradicate the tumour. This dilemma is embodied in the concept of therapeutic ratio, which represents the ratio of probabilities of tumour control versus normal tissue morbidity. Approaches to improving the therapeutic ratio include:                (a) optimising the physical targeting of the radiation to the tumour;        (b) fractionation of the radiation dose; and        (c) the use of radiomodifiers (which includes both radioprotectors and radiosensitisers, the latter of which can be used to increase the level of cell kill per unit of radiation dose).        \nImproving the physical delivery of radiation has had a considerable impact on the practice of radiotherapy. For example, increasing the energy of x-ray photons from several hundred kilovolts to the present-day megavoltage beams enables the zone of maximum radiation dose to be set at depths of several centimetres, whereas with the older machines the maximum dose was near the skin surface. There are a number of more sophisticated approaches to “tailoring” treatment beams in various stages of development and implementation. Brachytherapy, the use of implanted radioactive sources rather than external beams, is a further approach to improving the physical dose distribution.\nAlmost without exception, curative external beam radiotherapy involves fractionation of the radiation dose. An example of a conventional schedule would be a total of 60 Gray given in thirty 2 Gray fractions. Since cells have the capacity to repair radiation damage between fractions, the fractionated treatment results in much less cell-kill than a single dose of 60 Gray. However, normal cells generally have a greater repair capacity than do tumour cells, so the “sparing” effect of fractionation is more marked for normal tissues. In short, fractionation improves the therapeutic ratio.\nExploration of radiomodifiers such as radioprotectors and radiosensitisers has focussed on hypoxic cell sensitisers such as metranidazole and misonidazole. Radioprotectors have received much less attention than radiosensitisers at the clinical level. The nuclear era spawned considerable effort in the development of radioprotectors with more than 4000 compounds being synthesised and tested at the Walter Reed Army Institute of Research in the United States of America in the 1960's. With the exception of a compound that was called WR2728 (later called Ethyol and now known as Amifostine) none of the compounds have proved useful for cancer radiotherapy, and even WR2728 was considered too toxic for administration in either the military or industrial contexts (i.e., protection against total body irradiation). More recently, for example, Metz and co-workers (Metz et al, Clin Cancer Res. 10, 6411-17, 2004) (15) developed the radioprotective compound known as TEMPOL, which demonstrates only limited efficacy even at very high concentrations, and Burdelya and colleagues (Burdelya et al Science 320, 226-30, 2008) (16) developed the compound known as the TOLL receptor agonist which suffers from the necessity for it to be administered systemically.\nIt is important to note the interplay between the three approaches (a)-(c), above, to improving the therapeutic ratio. A combination of improved physical targeting, fractionation and radiomodifiers could transform the intent in some radiotherapy situations from palliative to curative. For curative schedules, successful application of radiomodifiers would relax the requirement for fractionation and hence reduce overall costs of treatment, which to a large extent is proportional to the number of treatment fractions per patient.\nA particularly important role for radioprotectors has emerged from the recognition that accelerated repopulation of tumour cells during radiotherapy can seriously compromise the effectiveness of treatment. The main consequences of this have been as follows:                (i) The development of accelerated treatment schedules to reduce the overall time of radiotherapy treatment. In such accelerated schedules, acute reactions are a particular problem. For example, acute oral mucositis in head and neck cancer patients indicates a clear need for radioprotectors.        (ii) The recognition that the interruption of radiotherapy treatment due to normal tissue reactions will reduce the probability of tumour control. Accordingly, the use of radioprotectors to prevent toxicity-induced treatment interruption would be clearly beneficial.        \nThe events of 11 Sep. 2001 prompted assessments of vulnerability to many types of terrorism scenarios, amongst which is a collection described as radiological terrorism. An example is the so-called “dirty bomb” involving dispersal of some form a radioactivity with conventional explosive. Whilst attention is focused on the acute radiation syndrome (ARS; also referred to as “radiation sickness”), which describes the consequences of whole-body exposure to radiation doses greater than 1 Gy, there are also concerns about the longer-term effects of low doses, namely radiation-induced mutagenesis and carcinogenesis (1). This general situation, and the realisation that no prophylactic agents are available to provide protection against exposure to ionising radiation, has generated significant research and political activity.\nThe mean lethal dose of radiation required to kill 50% of humans 60 days after whole-body irradiation (LD50/60) is between 3.25 and 4 Gy without supportive care, and 6-7 Gy when antibiotics and transfusion support are provided (1). The mortality is largely attributed to the haematopoietic syndrome, a consequence of hypoplasia or aplasia of the bone marrow. Cytopenias develop as a result of radiation-induced and normal attrition of mature functional cells, combined with the failure of replacement because of radiation-induced depletion of haematopoietic stem cells and progenitors. The time and extent of cytopenia generally correlate with radiation dose and prognosis, but the kinetics of depletion and recovery of blood cells also varies between the erythropoiesis, myelopoiesis and thrombopoiesis lineages, thrombopoiesis being the slowest.\nThe gastrointestinal syndrome results from ablation of stem cells in intestinal crypts, which in turn leads to denudation of the intestinal mucosa. This injury occurs after whole-body doses in the range of 3-15 Gy and in rodents doses at the upper end of this range usually result in death within about 1 week after irradiation.\nCountermeasures against unplanned irradiation include a wide range of potential molecular and cellular interventions. However, the mechanistic simplicity of chemical radioprotection—that is, reduction of radiation-induced DNA damage—is attractive because of its widespread potential. In this context, the possible need for protection of individuals at risk of exposure to low radiation doses, to thereby minimise long-term radiation effects such as mutagenesis and carcinogenesis, is particularly important. Such individuals would include emergency personnel involved in response to unplanned exposures, as well as those subject to occupational exposure to ionising radiation.\nA further group would be patients exposed to ionizing radiation during diagnostic medical procedures conducted in diagnostic radiology and nuclear medicine departments of hospitals and outpatient facilities.\nThe radioprotective properties of the minor groove binding DNA ligand Hoechst 33342 were first described by Smith, P. J. and Anderson, C. O. (2), who used clonogenic survival assays of irradiated cultured cells. Young, S. D. and Hill, R. P. (3) reported similar effects in cultured cells, but extended their studies to in vivo experiments. They concluded that the lack of radioprotection in their in vivo experiments was due to insufficient levels of Hoechst 33342 being delivered to target cells following intravenous injection. The findings of Hill and Young underline an important requirement for effective radioprotectors, namely potency. If the radioprotector is more potent, then it is more likely to achieve the required concentrations in an in vivo setting.\nThere is another aspect to be considered apart from potency. The concentration required for radioprotection must be non-toxic regardless of the potency of the radioprotector. If the radioprotector is delivered systemically, then this lack of toxicity requirement includes not just the cells and tissues to be protected from the radiation, but extends to the toxicity to the subject as a whole. In the case of Hoechst 33342 toxicity limits the extent to which it is useful as a radioprotector.\nThere is also a substantial conceptual problem in using radioprotectors in cancer radiotherapy. In attempting to decrease the effect of radiation on normal tissues by application of radioprotectors, there is a fear that some of the radioprotector will reach the tumour, thereby compromising tumour cell kill. The existing radioprotectors, e.g. WR2721 (Amifostine) and its active metabolite WR1065, are relatively small, diffusible molecules which do not avidly bind to tissue components and can therefore penetrate effectively through cell layers, so that they can reach the tumour via the circulation.\nThere is a need for radioprotectors that have limited penetration through cell layers. Such a property enables radioprotectors to be applied locally or topically to critical radiosensitive normal tissues in the vicinity of the tumour. Limited penetration restricts the extent to which the radioprotector reaches the capillary bed and is taken up into the circulation thereby reaching the tumour by systemic delivery in sufficient concentrations to confer significant radioprotection to the tumour.\nThe limited diffusion of DNA-binding ligands such as Hoechst 33342 through cell layers is known and has been exploited in mapping the location of cells in multi-cellular spheroids and in vivo, with respect to perfusion. Thus perfusion of Hoechst 33342 is considered a surrogate marker for perfusion of oxygen. In addition to restricting access to the tumour by systemic uptake following local or topical application to normal tissues, there is a further potential advantage of limited penetration in the context of cancer radiotherapy. This advantage stems from the view that the vasculature, in particular the endothelial cells, are the critical targets that determine the damaging effects of radiation. Furthermore, most radioresistant cells in the tumour are those viable cells that are most distant from the capillaries. The radioresistance of these cells is due to their hypoxic state, which in turn reflects their remoteness from the capillaries.\nConsequently, radioprotectors having limited diffusion, when administered intravenously, will be delivered more efficiently to critical radiosensitive cells in normal tissues, than to the hypoxic subpopulation of cells in tumours which limit the effectiveness of radiotherapy generally. Thus, the use of such radioprotectors would be expected to enable higher radiation doses to be used, with increased probability of killing the hypoxic cells in the tumour.\nHowever, the potential of the combination of these radiobiological features and the characteristics of DNA-binding radioprotectors can only be useful in cancer radiotherapy provided that an over-riding and necessary requirement of the radioprotectors exists, namely that the radioprotectors are sufficiently potent as to confer demonstrable radioprotection at non-toxic concentrations, when applied topically or systemically. A further practical requirement is that the extent of the limited penetration is sufficient to prevent significant systemic uptake following topical application, but not so pronounced so as to prevent sufficient concentrations from reaching the cells that determine the radiosensitivity of the tissue to be protected from the effects of ionising radiation, by topical or local application.\nThe extent of radioprotection (in the contexts of both cancer radiotherapy and protection from unplanned radiation exposure) is generally described in terms of dose modification factor (DMF), which is defined as the ratio of radiation doses required to produce the equivalent radiation-induced effect (molecular, cellular or in vivo endpoint) in the presence and absence of the radioprotector. When the radioprotective effect is observed on the basis of an in vivo endpoint, mechanisms other than modification of the initial radiation-induced damage may be involved. For example, for both the haematopoietic syndrome and the gastrointestinal syndrome, infection plays an important role in ultimate mortality, as a consequence of neutropenia and breach of the intestinal mucosal barrier, respectively. Thus, some immunostimulants have potential as mitigators of the radiation response. Immunostimulants can also be effective post-irradiation.\nInternational patent publication No. WO97/04776 and the subsequent publication by Martin et at (4) disclose certain bibenzimidazole compounds characterised by substitution with sterically hindering and electron donating groups. Although these compounds demonstrate strong radioprotective activity there is scope to reduce the inherent cytotoxicity of compounds of this general class. The challenge, however, is to do so while retaining, and preferably improving, radioprotective activity (measured as dose modification factor). The disclosures of WO97/04776 are included herein in their entirety by way of reference.\nInternational patent publication No. WO/2008/074091 also discloses bibenzimidazole compounds substituted with fluorine and/or chlorine and which, relative to known radioprotector compounds such as those of International patent publication No. WO97/04776, exhibit reduced cytotoxicity activity. While the cytotoxicity of the fluorine and chlorine substituted bibenzimidazole compounds was improved there is still a need for development of alternative radioprotective compounds, and preferably compounds that can be used in cancer radiotherapy, in protection of biological material from effects of radiation exposure and/or in protection of humans or animals from the effects of unplanned irradiation, which demonstrate low cytotoxicity but that retain radioprotective potency, and preferably that penetrate through cell layers to a limited extent. In particular it is desirable in some contexts that such compounds can be administered topically to protect tissues such as the skin, oral mucosa, oesophageal mucosa, rectal mucosa, vaginal mucosa and bladder epithelium, as well as parenterally to protect organs such as the lung and brain."}
{"text": "The present invention relates to a bisphenol-phosphorus compound complex and a thermally processed image recording material containing it.\nAs thermally processed image recording materials, there are a photothermographic material that has a photosensitive layer comprising a catalytically active amount of photocatalyst (e.g., silver halides), a heat developing agent, a reducible silver salt (e.g., silver salts of an organic acid), and a toning agent for controlling silver color tone as required, which are dispersed in a binder matrix, and a thermally processed image forming material that has a photosensitive layer comprising a heat developing agent, a reducible silver salt (e.g., silver salts of an organic acid), and a toning agent for controlling silver color tone as required, which are dispersed in a binder matrix. The thermally processed image forming material is heated imagewise at a high temperature to form monochromatic silver images through an oxidation-reduction reaction between the reducible silver salt (which functions as an oxidizing agent) and the heat developing agent. The photothermographic material is heated at a high temperature after imagewise light exposure to form monochromatic silver images through an oxidation-reduction reaction between the silver halide or the reducible silver salt (which functions as an oxidizing agent) and the heat developing agent. The oxidation-reduction reaction is accelerated by catalytic action of a latent image of silver halide generated upon exposure. Therefore, the monochromatic silver images are formed in exposed areas of the material. These materials are disclosed in many references including U.S. Pat. No. 2,910,377 and Japanese Patent Publication (Kokoku, hereinafter referred to as JP-B) 43-4924.\nSince the aforementioned thermally processed image recording materials are not subjected to a fixation treatment after the heat development, they suffer from a problem that the thermally reactive silver salt of an organic acid and reducing agent are left as they are in the thermally processed image recording materials, and thus white portions are colored when the materials are stored for a long period of time after the development. O-Bisphenol type reducing agents are frequently used for the thermally processed image recording materials, since they show high reactivity. If their amounts are reduced, the coloration of white portions can be effectively suppressed. However, if the amounts of the o-bisphenol type reducing agents are reduced, sufficient image density can no longer be obtained. Therefore, it is difficult to compromise the image storability and the image density.\nTherefore, an object of the present invention is to provide a thermally processed image recording material that can provide sufficient image density at a practical reaction temperature (specifically 100-140xc2x0 C.) and within a practical reaction time (specifically 1-30 seconds), and can sufficiently suppress the coloration of white portions when the material is stored in a dark place after the development. Further, another object of the present invention is to provide a compound useful as a reducing agent that can reconcile the image density and the image storability.\nThe inventors of the present invention assiduously studied in order to achieve the aforementioned objects. As a result, they found a surprising effect that if a complex obtained from a bisphenol compound, which is known as a reducing agent used for thermally processed image recording materials (see, for example, European Patent Publication EP0803764A1, Japanese Patent Laid-open Publication (Kokai, hereinafter referred to as JP-A) 51-51933, JP-A-6-3793 etc.), and a phosphorus compound is used, sufficient image density could be obtained and the image storability was markedly improved without substantially decreasing the reducing property. Thus, the present invention was accomplished.\nThat is, the present invention provides bisphenol-phosphorus compound complexes represented by the following general formula (1): \nIn the general formula (1), R1, R2, R3 and R4 each independently represent hydrogen atom or a group that can be a substituent on a benzene ring, and L represents xe2x80x94Sxe2x80x94 group or a xe2x80x94CHR5xe2x80x94 group. R5 represents hydrogen atom or an alkyl group. R6, R7 and R8 each independently represent an alkyl group, an aryl group, a heterocyclic group, xe2x80x94N(R9) (R10) or xe2x80x94O(R9). R9 and R10 each independently represent an alkyl group, an aryl group or a heterocyclic group. Two or more groups selected from R6, R7, R8, R9 and R10 may be bound together to form a ring.\nPreferred bisphenol-phosphorus compound complexes of the aforementioned general formula (1) are those wherein R1, R2, R3 and R4 each independently represent an alkyl group, L represents a xe2x80x94CHR5xe2x80x94 group, R5 represents hydrogen atom or an alkyl group, R6, R7 and R8 each independently represent an alkyl group, an aryl group or xe2x80x94O(R9), and R9 represents an alkyl group or an aryl group.\nThe present invention also provides a thermally processed image recording material comprising a non-photosensitive silver salt of an organic acid, a reducing agent for silver ions and a binder on one surface of a support, wherein the thermally processed image recording material comprises at least one kind of the bisphenol-phosphorus compound complex represented by the aforementioned general formula (I). The thermally processed image recording material of the present invention preferably contains at least one kind of photosensitive silver halide.\nThe thermally processed image recording material of the present invention utilizing the bisphenol-phosphorus compound complex represented by the following general formula (1) can provide sufficient image density at a practical reaction temperature (specifically 100-140xc2x0 C.) and within a practical reaction time (specifically 1-30 seconds), and can sufficiently suppress the coloration of white portions when the material is stored in a dark place after the development."}
{"text": "A cup of tea is typically made by dipping a porous tea bag into a hot cup of water for several seconds to several minutes in order to release the tea flavor from the tea bag or tea strainer into the hot water. The tea bag is then removed. Removing the tea bag can be a messy process as the tea bag tends to drip and inadvertently sink to the bottom of the cup. Moreover, drippings can scald the skin, stain the tablecloth and/or create a pool of tea in a saucer, which might then drip on and stain one's clothing or tablecloth. Other disadvantages with the common method of preparing a cup of tea include the length of time needed to achieve an optimal flavor and the singular and thus potentially wasteful use of tea bags. Similar issues arise when the cup of tea is prepared from tea leaves. To prevent the undesirable mess resulting from disposing of the used tea leaves, multiple dishes need to be typically used in a time consuming process. Moreover, the damp tea bag or tea strainer, once removed from the cup, is unattractive and detracts from the tea brewing experience.\nThe present invention addresses all these negatives. The device used for preparing the cup of tea from a tea bag is a single unit assembly that provides for functionalities and features to 1) prevent the tea bag or from sinking to the bottom so that its retrieval would not require reaching into the cup with ones' fingers or a spoon, 2) prevent dripping outside the cup when the bag is removed, 3) shorten the time to squeeze out the flavor from the tea bag into the water and obviate the need for wrapping the tea bag around a spoon or other secondary utensil to squeeze the tea bag before removing from the cup of tea, 4) enable reusing a tea bag at least once, and 5) obscure from view the unattractive, damp, used tea bag following brewing.\nThe device for preparing the cup of tea from tea leaves is a single unit assembly that provides for functionalities and features to 1) prevent the tea leaves from sinking to the bottom of the cup so that its retrieval would not require reaching into the cup with a spoon, 2) prevent dripping outside the cup when the bag is removed, and 3) shorten the time to squeeze out the flavor from the tea leaves into the water and obviate the need for handling the tea leaves multiple times. 4) potentially reuse the tea leaves at least once and 5) obscure from view the unattractive, damp, used tea leaves following brewing.\nA number of prior art references that relate to tea bag holding and squeezing devices. In U.S. Pat. No. 5,335,591 two porous plates are attached to the ends of the tong. U.S. Pat. No. 2,887,948 likewise discloses two dome shaped porous plates connected to a common hinge. References that indicate at least one porous plate include U.S. Pat. No. 3,342,518 and U.S. Pat. No. 2,708,401."}
{"text": "The present invention relates to a charge-pump circuit outputting voltage fluctuation within a step or power source voltage Vdd and a control method thereof, particularly a control method of a charge-pump circuit capable of normal charging pump operation removing influence of a parasitic diode following to a charge transfer device.\nThe charge-pump circuit developed by Dicson generates higher voltage than power source voltage Vdd of an LSI chip by voltage fluctuation of each pumping packet connecting plural stages of the pumping packet in series. For example, it is used for generating voltage for program/erase of flash memories.\nHowever, the conventional charge-pump circuit carries out voltage fluctuation with a step of the power source voltage Vdd, and a circuit capable of carrying out voltage fluctuation with lower voltage step than the Vdd was not yet proposed. So, the inventor has already proposed a charge-pump circuit capable of carrying out voltage fluctuation with lower voltage step than the Vdd and improving efficiency xcex7 of the circuit (U.S. patent application No. 09/732,944 filed on Dec. 8, 2000).\nThe outline thereof will be described below. FIG. 10 to FIG. 12 are circuit diagrams showing a structure of an operation of xe2x88x920.5 Vdd voltage fluctuation charge-pump circuit. The charge-pump circuit generates voltage of voltage fluctuation of xe2x88x920.5 Vdd to earth voltage (0 V).\nIn FIG. 10, diodes D1 and D2 are connected in series as charge transfer devices. To a cathode of the diode D1, earth voltage (0 V) is supplied. The diodes D1 and D2 generally consist of MOS transistors for charge transfer in order to integrate into an LSI.\nSwitches S1, S2 and S3 connect two capacitors 1 and 2 to a connecting point between the diodes D1 and D2 switching in parallel or in series. These switches S1, S2 and S3 can consist of MOS transistors. Thus, on and off of the switches S1, S2 and S3 corresponds to on and off of the MOS transistors. A clock driver 3 supplies clock CLK to the capacitor 2. Output voltage output from the diode 2 is applied to a load 4.\nAn outline of control method of the charge-pump circuit will be described below. Power source voltage of the clock driver is assumed 5 V. Although forming the diodes D1 and D2 and the switches S1, S2 and S3 actually occurs voltage drop, the voltage drop is assumed 0 V omitting the voltage drop here.\nWhen input clock of the clock driver 3 is high level (CLK=High), assuming that S1 is off, S2 is on and S3 is off, two capacitors 1 and 2 are connected in series and each node voltage is: VL1≈0V, VA=VB=2.5V, VC=5V. VL1 is voltage of a connecting node between the diode D1 and the capacitor 1 (a pumping node), VA is voltage of a connecting node between the capacitor C1 and the switch S2, VB is voltage of a connecting node between the switch S2 and the capacitor 2, and VC is voltage of a connecting node between output of the clock driver 3 and the capacitor 2.\nThat is, if capacitance values which capacitors 1 and 2 have are equal, capacitors 1 and 2 are respectively charged to voltage of Vdd/2 by distributing equally electric charge to the capacitors 1 and 2 (see FIG. 10).\nNext, when S2 is off and S1 and S3 are on in the state of CLK=High, two capacitors 1 and 2 are switched to parallel connection. Thus, each node voltage becomes: VL1≈2.5V, VA=5V, VB=2.5V, VC=5V (See FIG. 11).\nNext, when the input clock CLK is transferred to low level (CLK=Low) in the state of the parallel connection, each node voltage becomes: VL1≈2.5V, VA=0V, VB=xe2x88x922.5V, VC=5V by effect of the capacitor coupling because the capacitors 1 and 2 are connected to the pumping node (See FIG. 12).\nThus, by repeating switching the capacitors 1 and 2 alternately to series and parallel according to the input clock CLK, output voltage of xe2x88x922.5 V (=(xe2x88x92xc2xd)xc2x7Vdd) is supplied to the load 4 from the diode D2.\nWhen the diodes D1 and D2 consist of MOS transistors for charge transfer where a source and a gate are connected, there are a problem that needless current flows transitionally by that the diode D1 is biased to forward direction when voltage of the pumping node VL1 becomes 2.5V. Then, in order to avoid the problem, gate voltage of the MOS transistor for charge transfer may be controlled separating from source voltage.\nAt timing connecting the capacitors 1 and 2 in series, gate voltage of the transistor for charge transfer suitable for the diode D1 is made on by setting low level (see FIG. 10), at timing connecting the capacitors 1 and 2 in parallel, gate voltage of the transistor for charge transfer suitable for the diode D1 is made off by setting high level (see FIG. 11).\nHowever, in the above-mentioned control method of the charge-pump circuit, voltage of the pumping node VL1 repeats change such as 0Vxe2x86x922.5Vxe2x86x92xe2x88x922.5V. Because of that, even if the MOS transistors for charge transfer are any of P-channel and N-channel, a problem occurs that a parasitic diode formed incidentally to the MOS transistors is biased to forward direction and voltage fluctuation is not carried out normally.\nFIGS. 13A and 13B are views showing a problem in the case that the diode D1 is made by P-channel MOS transistor as a charge transfer device. In this case, a source S and a substrate B are earthed to improve efficiency of the charge-pump circuit depressing back gate bias effect of the MOS transistor.\nAs shown in FIG. 13A, there is not any problem in case that voltage of the pumping node VL1 is xe2x88x922.5V. However, as shown in FIG. 13B, when a parasitic diode formed between a drain and the substrate in the case that voltage VL1 is 2.5V of the pumping node is biased to forward, forward direction current of the diode flows between the drain and the substrate, power efficiency becomes bad, and charge-pump operation is not carried out normally.\nFIGS. 14A and 14B are views showing a problem in the case that the diode D1 is made by N-channel MOS transistor as a charge transfer device. In this case, a drain D (pumping node) and a substrate B are connected to depress back gate bias effect of the MOS transistor.\nAs shown in FIG. 14A, there is not any problem in case that voltage of the pumping node VL1 is xe2x88x922.5V. However, as shown in FIG. 14B, a parasitic diode formed between the substrate and a source is biased to forward in the case that voltage VL1 is 2.5V. Then, forward direction current of the diode flows between the drain and the substrate, power efficiency becomes bad, and charge-pump operation is not carried out normally.\nAn object of the invention is to prevent that the parasitic diode is biased to forward direction and needless current flows, and to make normal operation of the charge-pump circuit possible.\nA charge-pump circuit of the invention comprises, at least first and second MOS transistors for charge transfer connected in series, first and second capacitors, clock supplying means supplying clock to one end of the second capacitors, first switching means for connecting said first and second capacitors to a connecting point of the first and second MOS transistors for charge transfer in series, and second switching means for connecting said first and second capacitors to the connecting point of the first and second MOS transistors for charge transfer in parallel, wherein said clock supplying means changes the state of said clock when said first and second switching means turn off.\nBy such the structure, a timing that the clock is supplied to the capacitor changes to high level from low level (or to low level from high level) is adjusted in the state that both of the first and second switch means are off. In this state, the first and second capacitors are separated from the connecting point (pumping node) of the first and second MOS transistors for charge transfer.\nThus, it is prevented that the parasitic diodes following the first and second MOS transistors for charge transfer are biased to forward direction because change of potential of the pumping node is depressed."}
{"text": "Uroguanylin, guanylin and bacterial ST peptides are structurally related peptides that bind to a guanylate cyclase receptor and stimulate intracellular production of cyclic guanosine monophosphate (cGMP) (1-6). This results in the activation of the cystic fibrosis transmembrane conductance regulator (CFTR), an apical membrane channel for efflux of chloride from enterocytes lining the intestinal tract (1-6). Activation of CFTR and the subsequent enhancement of transepithelial secretion of chloride leads to stimulation of sodium and water secretion into the intestinal lumen. Therefore, by serving as paracrine regulators of CFTR activity, cGMP receptor agonists regulate fluid and electrolyte transport in the GI tract (1-6; U.S. Pat. No. 5,489,670).\nThe process of epithelial renewal involves the proliferation, migration, differentiation, senescence, and eventual loss of GI cells in the lumen (7,8). The GI mucosa can be divided into three distinct zones based on the proliferation index of epithelial cells. One of these zones, the proliferative zone, consists of undifferentiated stem cells responsible for providing a constant source of new cells. The stem cells migrate upward toward the lumen to which they are extruded. As they migrate, the cells lose their capacity to divide and become differentiated for carrying out specialized functions of the GI mucosa (9). Renewal of GI mucosa is very rapid with complete turnover occurring within a 24-48 hour period (9). During this process mutated and unwanted cells are replenished with new cells. Hence, homeostasis of the GI mucosa is regulated by continual maintenance of the balance between proliferation and apoptotic rates (8).\nThe rates of cell proliferation and apoptosis in the gut epithelium can be increased or decreased in a wide variety of different circumstances, e.g., in response to physiological stimuli such as aging, inflammatory signals, hormones, peptides, growth factors, chemicals and dietary habits. In addition, an enhanced proliferation rate is frequently associated with a reduction in turnover time and an expansion of the proliferative zone (10). The proliferation index has been observed to be much higher in pathological cases of ulcerative colitis and other GI disorders (11). Thus, intestinal hyperplasia is the major promoter of gastrointestinal inflammation and carcinogenesis.\nIn addition to a role for uroguanylin and guanylin as modulators of intestinal fluid and ion secretion, these peptides may also be involved in the continual renewal of GI mucosa. Previously published data in WO 01/25266 suggests a peptide with the active domain of uroguanylin may function as an inhibitor of polyp development in the colon and may constitute a treatment of colon cancer. However, the mechanism by which this is claimed to occur is questionable in that WO 01/25266 teaches uroguanylin agonist peptides that bind specifically to a guanylate cyclase receptor, termed GC-C, that was first described as the receptor for E. coli heat-stable enterotoxin (ST) (4). Knockout mice lacking this guanylate cyclase receptor show resistance to ST in intestine, but effects of uroguanylin and ST are not disturbed in the kidney in vivo (3). These results were further supported by the fact that membrane depolarization induced by guanylin was blocked by genistein, a tyrosine kinase inhibitor, whereas hyperpolarization induced by uroguanylin was not effected (12,13). Taken together these data suggest that uroguanylin also binds to a currently unknown receptor, which is distinct from GC-C.\nOther papers have reported that production of uroguanylin and guanylin is dramatically decreased in pre-cancerous colon polyps and tumor tissues (14-17). In addition, genes for both uroguanylin and guanylin have been shown to be localized to regions of the genome frequently associated with loss of heterozygosity in human colon carcinoma (18-20). Taken together, these findings indicate that uroguanylin, guanylin and other peptides with similar activity may be used in the prevention or treatment of abnormal colon growths. This proposal is bolstered by a recent study demonstrating oral administration of uroguanylin inhibits polyp formation in mice (15,16).\nUroguanylin and guanylin peptides also appear to promote apoptosis by controlling cellular ion flux. Alterations in apoptosis have been associated with tumor progression to the metastatic phenotype. While a primary gastrointestinal (GI) cancer is limited to the small intestine, colon, and rectum, it may metastasize and spread to such localities as bone, lymph nodes, liver, lung, peritoneum, ovaries, brain. By enhancing the efflux of K+ and influx of Ca++, uroguanylin and related peptides may promote the death of transformed cells and thereby inhibit metastasis.\nOne of the clinical manifestations of reduced CFTR activity is the inflammation of airway passages (21). This effect may be due to CTFR regulating the expression of NF-KB, chemokines and cytokines (22-25). Recent reports have also suggested that the CFTR channel is involved in the transport and maintenance of reduced glutathione, an antioxidant that plays an important role in protecting against inflammation caused by oxidative stress (39). Enhancement of intracellular levels of cGMP by way of guanylate cyclase activation or by way of inhibition of cGMP-specific phosphodiesterase would be expected to down-regulate these inflammatory stimuli. Thus, uroguanylin-type agonists should be useful in the prevention and treatment of inflammatory diseases of the lung (e.g., asthma), bowel (e.g., ulcerative colitis and Crohn's disease), pancreas and other organs.\nOverall, it may be concluded that agonists of guanylate cyclase receptor such as uroguanylin have potential therapeutic value in the treatment of a wide variety of inflammatory conditions, cancer (particularly colon cancer) and as anti-metastatic agents. The development of new agonists is therefore of substantial clinical importance."}
{"text": "This invention relates in general to bulk material handling devices and in particular to a new and useful blowing device for eliminating compactions in storage silos for bulk materials.\nThe invention concerns a blowing device which uses air blasts for the elimination of compactions in bulk material storage silos with a compressed air storage container, with a compressed air supply duct connected to the compressed air storage container, a compressed air outflow duct, a rapid exit valve interposed between the compressed air supply duct and the compressed air outflow duct and with a compressed air inlet pipe socket with a multiway valve connected to the rapid exit valve whereby in the rapid exit valve between open position and closed position a directly air pressure actuated valve closure is guided and in a closed position, is pressed against the compressed air outflow duct and is fashioned as a valve seat.\nA blowing device is known of a kind in which the compressed air supply duct, the compressed air outflow duct and the rapid exit valve and its housing presents a construction unit. In this construction unit the compressed air supply duct and the compressed air outflow duct are arranged orthogonally to each in such a way that during afflux of the rapid exit valve considerable flow and throttle losses occur. The same is true of the form and guidance of the valve closure which is practically fashioned as a valve disk with an inside and an outside collar. In a closed position the valve disk contacts, with the inside collar, the valve seat of the compressed air outflow duct at the end side while the outside collar is intended to guide the valve disk in the valve housing. This guiding, however, is already somewhat problematic if for no other reason than between the outside collar and the inner wall of the valve housing, sufficient margin clearance must be given to allow motion so that the compressed air storage container can be supplied with compressed air through the compressed air inlet pipe socket connected to the valve housing with the rapid exit valve in closed position. An annular channel between the outside collar of the valve disk and the guiding surface on the inside wall of the valve housing practically serves this purpose. The resulting motion can lead to flutter effects or tilting in the case of valve disks and consequently to reduce flow velocities in the compressed air outflow duct and to greatly slowed or completely missing opening processes such that an impulse-like discharge of the compressed air storage container is frequently hardly achieved. However, air blasts then remain to a large extent ineffective. In this respect the invention proposes a solution."}
{"text": "1. Technical Field\nThe present disclosure relates to an operational amplifier and, more particularly, to a low-voltage operational amplifier and an operational amplifying method.\n2. Discussion of Related Art\nAs techniques in a semiconductor manufacturing process are improved, the width between the semiconductor circuit lines is gradually decreased. Because the width between the semiconductor circuit lines becomes narrower, an integration degree of the semiconductor is improved and a production cost of the semiconductor is reduced. In addition, as the width between the semiconductor circuit lines gradually becomes narrower, a power-supply voltage that is applied to the semiconductor gradually becomes lower. For example, a semiconductor that has a 90 nm width between the semiconductor circuit lines generally needs a power-supply voltage that is not more than 1.2V. On the other hand, a semiconductor that has a 90 nm width between the semiconductor circuit lines occasionally needs a power-supply voltage that is not more than 1.0V. Therefore, a circuit that operates properly at a high power-supply voltage may not operate properly at a low power-supply voltage. Accordingly, a circuit that has a new structure is needed.\nFIG. 1 is a diagram illustrating a conventional operational amplifier that has a folded cascode structure.\nThe conventional operational amplifier usually includes a multi-stage structure in order to achieve a high gain. The conventional operation amplifier shown in FIG. 1 includes two stages 10 and 20 that each have a folded cascode structure.\nThe conventional operational amplifier of FIG. 1 achieves the high gain owing to the cascade connected stage. There are two cascade connected stages in each stage 10 and 20, and one such cascade stage 30 has four MOS transistors connected in series between a power-supply voltage VDD and a ground voltage VSS. A second cascade connected stage of four MOS transistors is in the first stage 10 and bias voltages BS1, BS2, and BS3 are fed to respective gates of corresponding transistors and a common mode feedback signal is connected to the gates of corresponding transistors whose drains are connected to the ground voltage VSS. The input signals IN+ and IN− are fed to an input pair of transistors that have their drains connected to a current source transistor that receives at its gate a bias voltage BS4. Outputs of the first stage 10 are fed to a pair of input transistors whose drains are commonly connected to a current source transistor whose gate receives the bias voltage BS4 and whose drain is connected to the ground voltage VSS. The second stage 20 employs the same structure as the first stage 10 with corresponding transistors receiving the bias voltages BS1, BS2, and BS3 and the common mode feedback signal CMFB. Outputs 40 the second stage 20 are at the same location as the outputs of the first stage 10, with the outputs OUT+ and OUT− of the first stage being coupled with respective capacitors. As is known, a common mode feedback signal decides a common mode voltage of an amplified first signal and an amplified second signal. Thus, the conventional operational amplifier of FIG. 1 does not efficiently obtain a large dynamic voltage margin at a relatively low power-supply voltage. In addition, a swing margin at the output terminal 40 of the second cascade stage 20 is limited at a low power-supply voltage.\nFIG. 2 is a diagram illustrating another conventional operational amplifier.\nThe conventional operational amplifier of FIG. 2 has a structure where a folded cascode stage 110 is connected to a common source stage 120. Therefore, the conventional operational amplifier of FIG. 2 can obtain a larger swing margin at an output terminal than the conventional operational amplifier of FIG. 1. In a folded cascode stage 110, four MOS transistors 130 are connected in series between a power-supply voltage VDD and a ground voltage VSS, as in FIG. 1. A second set of four MOS transistors are provided in the first stage 110 and respective gates of corresponding transistors receive bias voltages BS1, BS2, and BS3, as well as the common mode feedback signal CMFB. The input signals IN+ and IN− are fed to a transistor pair whose common drains are connected to a current source transistor that receives at a gate thereof a bias voltage BS4 and that has a drain connected to the ground voltage VSS. The common source stage 120 has a pair of input MOS transistors having respective gates receiving the output signals from the folded cascode stage 130. A transistor pair has its drains connected to the power-supply voltage VDD, sources connected to the input transistor pair, and gates receiving the bias voltage BS1. A current source transistor is connected between the common drains of the input transistor pair and the ground voltage VSS with the common mode feedback signal CMFB connected to the gate thereof, the outputs OUT+ and OUT− of the operational amplifier are taken from the common source stage 120 by being capacitor coupled from the outputs of the folded cascode stage 110. Thus, the conventional operational amplifier of FIG. 2 does not efficiently obtain a large dynamic voltage margin at a low power-supply voltage. In addition, the conventional operational amplifier of FIG. 2 achieves a relatively low gain compared with the other conventional operational amplifier of FIG. 1.\nAccordingly, a low-voltage operational amplifier that can operate properly at a low power-supply voltage and that can have a high gain is needed."}
{"text": "Some methods have been developed with a view to obtaining a quasi-isothermal compression or expansion process, according to which, in order to obtain the theoretical condition of an isothermal transformation, i.e., the maintenance of equality between the mechanic work received during the compression phase or yielded during the expansion phase and the heat evacuated during the compression phase or the heat absorbed during the expansion phase respectively, the work space of variable size of an engine has been connected to a cooled heat exchanger, consisting of one or more heat exchange units, in series, during the compression phase and a heated heat exchanger during the expansion phase (U.S. Pat. No. 3,867,815). This method has the disadvantage that the volume of the heat exchangers adds to the volume of the dead space, detemined by the constructive parameters of the work space of variable size, thus preventing high compression ratios from being reached. In addition, owing to the fact that only one heat-exchanger is used, the equality between the received or transferred machanic work and the evacuated or absorbed heat respectively, cannot be ensured at any instant, consequently, the transformation curve moves significantly away from the theoretical isothermal curve, thereby damaging the efficiency of the cycle on the whole. Then there are also Stirling external combustion engines built according to different principles, in which, after the compression phase, the working agent is cooled inside a heat exchanger, afterwords run through a regenerator and finally introduced into a heated expansion space (Stirling engine, by G. Walker). This type of external combustion engines has the disadvantage of not being able to reach higher compression values, thereby affecting the general output of the engine."}
{"text": "1. Field\nThe present application relates to a medical handpiece, in particular a dental handpiece, having a lighting device comprising at least one optical semiconductor element.\n2. Description of Prior Art\nJapanese Patent Application JP 10165419 A describes a dental handpiece having a lighting device capable of emitting radiation of different wavelengths to thereby be able to perform different treatments. The handpiece therefore comprises multiple light sources in the form of LEDs and a switch by means of which optionally one or more LEDs can be operated so that green or blue light, for example, is emitted for treatments or white light is emitted for illuminating the treatment site.\nOne disadvantage is that the design of this handpiece is complicated and requires a great deal of space because multiple LEDs, a switch and electronic or mechanical switch elements have to be accommodated in the handpiece. Because of a wide variety of components contained in the handpieces with a connection for tools that can be set in motion, e.g., drive elements, tool receptacles, media lines, etc., integration of additional elements is not possible to an unlimited extent or would be possible only if the outside dimensions of the handpiece could be enlarged. However, the space available in the oral cavity is very limited, especially for dental applications, and the user's vision is often limited to the preparation site, so larger outside dimensions are not desirable.\nOne object is thus to create a medical handpiece, in particular a dental handpiece, having a receptacle for tools that can be set in motion and having a lighting device with which treatments can be performed under different operating parameters, e.g., different wavelengths, whereby the handpiece should have the simplest possible design and there should not be any increase in the outside dimensions of the handpiece."}
{"text": "Duplex stainless steel refers to ferritic austenitic steel alloy. Such steels have a microstructure comprising ferritic and austenitic phases. The duplex steel alloy, to which the invention pertains, is characterized by a high content of Cr and N and a low content of Ni. Background references in this respect include WO 95/00674 and U.S. Pat. No. 7,347,903. The duplex steels described therein are highly corrosion resistant and can therefore be used, e.g., in the highly corrosive environment of a urea manufacturing plant.\nUrea (NH2CONH2) can be produced from ammonia and carbon dioxide at elevated temperature (typically between 150° C. and 250° C.) and pressure (typically between 12 and 40 MPa) in the urea synthesis section of a urea plant. In this synthesis, two consecutive reaction steps can be considered to take place. In the first step, ammonium carbamate is formed, and in the next step, this ammonium carbamate is dehydrated so as to provide urea, The first step (i) is exothermic, and the second step can be represented as an endothermic equilibrium reaction (ii):2NH3+CO2→H2N—CO—ONH4  (i)H2N—CO—ONH4H2N—CO—NH2+H2O  (ii)\nIn a typical urea production plant, the foregoing reactions are conducted in a urea synthesis section so as to result in an aqueous solution comprising urea. In one or more subsequent concentration sections, this solution is concentrated to eventually yield urea in the form of a melt rather than a solution. This melt is further subjected to one or more finishing steps, such as prilling, granulation, pelletizing or compacting.\nA frequently used process for the preparation of urea according to a stripping process is the carbon dioxide stripping process, as for example described in Ullmann's Encyclopedia of Industrial Chemistry, Vol. A27, 1996, pp 333-350. In this process, the synthesis section is followed by one or more recovery sections. The synthesis section comprises a reactor, a stripper, a condenser and, preferably but not necessarily, a scrubber in which the operating pressure is in between 12 and 18 MPa, such as in between 13 and 16 MPa. In the synthesis section, the urea solution leaving the urea reactor is fed to a stripper in which a large amount of non-converted ammonia and carbon dioxide is separated from the aqueous urea solution.\nSuch a stripper can be a shell- and tube-heat exchanger in which the urea solution is fed to the top part at the tube side and a carbon dioxide feed, for use in urea synthesis, is added to the bottom part of the stripper. At the shell side, steam is added to heat the solution. The urea solution leaves the heat exchanger at the bottom part, while the vapor phase leaves the stripper at the top part. The vapor leaving said stripper contains ammonia, carbon dioxide, inert gases and a small amount of water.\nSaid vapor is condensed in a falling film type heat exchanger or a submerged type of condenser that can be a horizontal type or a vertical type. A horizontal type submerged heat exchanger is described in Ullmann's Encyclopedia of Industrial Chemistry, Vol. A27, 1996, pp 333-350. The formed solution, which contains condensed ammonia, carbon dioxide, water and urea, is recirculated together with the non-condensed ammonia, carbon dioxide and inert vapor.\nThe processing conditions are highly corrosive, particularly due to the hot carbamate solution. In the past, this presented a problem in the sense that the urea manufacturing equipment, even though made from stainless steel, would corrode and be prone to early replacement.\nThis has been resolved, particularly by making the equipment, i.e. the relevant parts thereof subjected to the mentioned corrosive conditions, from a duplex steel described in WO 95/00674 (also known by the trademark of Safurex®). However, even though the foregoing reflects a major advancement in urea production, a particular problem exists in the stripper. A typical carbamate stripper comprises a plurality (several thousand) of tubes. Through the tubes, a liquid film runs downwards whilst stripping gas (typically CO2) runs upwards. Provisions are generally made to ensure that all tubes have the same load of liquid so as to have a flow of the liquid at the same speed. For, if the liquid does not flow through all of the tubes at the same speed, the efficiency of the stripper is reduced. These provisions comprise a liquid distributor, generally in the form of a cylinder with small holes in it.\nIt has been experienced that the liquid distributors need a relatively frequent replacement. Particularly, the size and shape of the holes changes with time, apparently as a result of corrosion, despite the fact that the liquid distributors are made from corrosion-resistant duplex steel as mentioned above. Thus, the affected distributors result in a different throughput of liquid in the stripper, as a result of which the desired equal loading of the stripper's tubes is less efficient.\nIt is therefore desired in the art to provide a corrosion resistant material that would provide the liquid distributors in the stripper with a better corrosion endurance."}
{"text": "The present invention relates generally to the packaging of integrated circuits. More particularly, the invention relates to dice having integrated pedestals.\nThere are a number of conventional processes for packaging integrated circuits. In many situations it is desirable to incorporate multiple integrated circuit dice into the same package in what is commonly referred to as a multi-chip package. Some multi-chip packages are arranged to stack two or more dice on top of each other. These stacked die packages have several potential advantages including the possibility of a reduced die or package footprint and certain performance advantages. For example, reducing the path length of electrical connections between integrated circuits potentially increases speed and reduces inductance of inter-chip communications.\nOne existing stacked die configuration is generally illustrated in FIG. 1. In this configuration a first die 104 is mounted on a carrier such as a lead frame 106 or a planar substrate. A second die 105 is then adhesively secured to the top surface of the first die 104 thereby creating a stacked die configuration. Bonding wires 108 are used to electrically connect both of the dice 104, 105 to the carrier (lead frame 106) using conventional wire bonding. An encapsulant material 115 is molded over the stacked dice to create a stacked die package 120. This stacked die approach allows the top die 105 to be electrically connected to both the underlying base die 104 and to a substrate or lead frame 106. This arrangement also generally requires that the top die 105 must be smaller than the base die 104 so that the top die 105 does not cover any of the bond pads on the base die 104.\nAnother stacked die approach is generally illustrated in FIG. 2. In this approach, a spacer 130 is adhesively secured to (or dispensed onto) the top surface of the bottom die 104. The top die 105 is then adhesively adhered to the spacer 130. The spacer 130 is sized so that it extends higher than the loop height of the bonding wires 108 used to electrically connect the base die, thereby providing clearance for the bonding wires. The spacer 130 can be formed from virtually any material that is compatible with the dice. By way of example, the spacer may be formed from a variety of materials such as silicon, various ceramics, etc. Alternatively, the spacer may be formed in situ by dispensing an adhesive material having ball like spherical support structures therein on the top surface of the lower die 104 in a region that is interior to the bond pads. The spacer 130 effectively forms a plateau on the bottom die that provides structural support for the top die. With this arrangement, the top die 105 can be virtually any size including the same size or even larger than the base die.\nAlthough the foregoing techniques work well in many applications, in the semiconductor industry, there are continuing efforts to provide more efficient packaging approaches. The packaging arrangements described below are particularly useful in forming stacked multi-chip packages."}
{"text": "A network camera system is known, which displays an image captured by a camera on a terminal apparatus such as a personal computer in real time and controls the image capturing direction of the camera from the terminal apparatus by using the Internet technology. There exists a monitoring system which records and reproduces a camera image from a remote site by using the network camera system.\nTo display a plurality of recorded images on a single window, the conventional monitoring system employs a method of, e.g., simultaneously recording/reproducing the contents of a plurality of images on a magnetic tape.\nHowever, if, for example, a user wants to reproduce a camera image of a specific time while simultaneously reproducing another camera image of a specific time and compare the images, the operation is complex."}
{"text": "Systems and devices often rely upon two kinds of signals for operation. A power signal may be used to provide power, logic, and control signals within a system or device. Further, a timing signal or clock signal may also be used to trigger, manipulate or otherwise control various components and circuits of a system or device. Thus, together, these two kinds of signals may be used accordingly.\nOften times, a system or device may have the above-described signals provided externally from the system or device. That is, the power signal and/or the clock signal may originate from a circuit outside the system or device. For example, when monitoring overhead power lines, a device for monitoring may draw power directly from the power lines themselves. Similarly, in a device having multiple integrated circuit (IC) chips or multiple separate printed circuit boards (PCB), a single clock source may provide timing signals for all components in a device. Furthermore, these signals, when originating outside of ICs, PCBs, or specific separate devices, may require isolation before being used within these components. As such, both the power signals and the timing signals may require isolating transformers for each kind of signal to provide the isolated signals for use on-chip. Providing two isolating transformers for these two sets of signals is cumbersome and inefficient."}
{"text": "Needle roller bearings are often used for planetary gear bearings of a planetary gear drive in the consideration of capacity, shape and the like. Recently, full roller type needle roller bearings are often used to increase load capacity.\nPatent document 1 proposes a lubrication structure for needle roller bearings used as the planetary gear bearing of a planetary gear drive. In this document, it is described that spiral lubricant retaining grooves are formed on outer peripheral surfaces of needle rollers.\nPatent Document 1: JP 2007-71355 A\nWhen needle roller bearings, especially, those of full roller type are used for a planetary gear bearing, lubricant retaining space is limited, so that there is a risk of poor lubrication if grease lubrication is employed."}
{"text": "The present invention relates generally to optical probes and more particularly to optical probes for use with Optical Coherence Tomography (OCT) and other optical imaging modalities.\nOCT is a laser based imaging modality that uses near infrared or infrared laser light to non-destructively image subsurface tissue structures. An imaging depth on the order of millimeters (mm), with a spatial resolution of a few micrometers (μm) is relatively easily achieved using OCT at practical light fluence levels on the order of 100 μW. OCT is therefore very useful for in vitro and in vivo tissue structure imaging applications such as may be used during minimally invasive surgical procedures. Currently, both side-imaging endoscope systems and forward imaging endoscope systems are known.\nThe construction of a needle endoscope that is capable of performing forward OCT imaging presents very significant design challenges. Current endoscopes are typically more than 5 mm thick. The thickness of such probes, especially when compared with their en face imaging area, e.g., about 2 mm wide, makes them undesirable as a needle endoscope for image-guided surgical procedures. One major challenge of making a thin endoscope lies with the difficulty of designing a probe beam deflection system that is capable of covering a sufficient scan volume while constraining the probe diameter to be less than about 2 mm to minimize the invasiveness of the probe. A reasonable OCT scan volume for providing sufficient image information would be a conical volume that is about 3 mm in length and about 2 mm in diameter at its maximum circumference.\nTherefore it is desirable to provide probes such as forward imaging endoscope needles useful for OCT imaging of a scan volume that overcome the above and other problems."}
{"text": "The present invention relates to the art of petroleum processing and, in particular, to the disposal of harmful and noxious waste products resulting therefrom.\nCrude oils are exceedingly complex mixtures, consisting predominantly of hydrocarbons containing sulfur, nitrogen, oxygen, and metals as minor constituents. While it is desirable to recover the hydrocarbon constituents in their pure form, realistically it is very difficult to isolate pure products because most of the minor constituents occur in combination with carbon and hydrogen. Separation of impurities such as those listed above generally requires expenditures of valuable resources such as time, chemicals, energy, and money. Therefore, it is the constant goal of the petroleum processing industry to optimize impurity-removal procedures, equipment, and resources in order to eliminate those impurities which have the most degrading effect on the end products.\nPerhaps the most ubiquitous impurity encountered in petroleum processing is sulfur. The presence of sulfur in petroleum products and, indeed, in the crude feedstock itself generally increases the corrosive characteristics thereof, and forms harmful and noxious reaction products upon combustion. In particular, the presence of sulfur-containing compounds reduces the combustion characteristics of gasoline and may render fuel oil unusuable in many places due to local regulation on the amount of sulfur allowed therein. Consequently, at nearly every stage of production measures are taken to either reduce the amount of sulfur or to render the sulfurcontaining compounds inoffensive.\nOne method for removing sulfur-containing compounds--hydrogen treating of petroleum fractions--has been known since the 1930's. However, it was not until the advent of catalytic reforming, which made inexpensive hydrogen-rich off-gas available, that hydrogen desulfurization developed to commercial level. Presently, hydrogen desulfurization is primarily associated with a catalytic reaction using colbalt molybdate on an alumina carrier. The feedstock is mixed with recycle and make-up hydrogen and heated to 400.degree.-850.degree. F., then charged to a fixed bed reactor at 50-1,500 psig.\nHydrogen treating is now used extensively to prepare reformer feedstock and, to some extent, for catalytic cracking feedstock preparation. It may also be used to upgrade middle distillates, cracked fractions, lube oils, gasolines, and waxes. Hydrodesulfurization, however, is a high energyconsuming process which also requires a supply of hydrogen.\nMoreover, a major effluent resulting from hydrodesulfurization is hydrogen sulfide, H.sub.2 S,--a flammable poisonous gas. Even though hydrogen sulfide may simply be burned-off into the atmosphere legislation in recent years has effectively limited this method of disposal because of the formation of sulfur dioxide which is intensely irritating to the eyes and respiratory system. Accordingly, alternative means for disposing of hydrogen sulfide have been developed and implemented.\nThe primary method of disposing of hydrogen sulfide is to convert the sulfur-bearing gas to elemental sulfur and water by, for instance, the Claus process. While this alternative may appear somewhat attractive since elemental sulfur is a saleable commodity, the Claus process requires construction of sulfur plants, quantities of catalysts, and energy. Furthermore, the market for elemental sulfur is not so extensive as to be able to absorb all the elemental sulfur currently produced without depressing the price therefor. Finally, the Claus process itself is fraught with some difficulty in that approximately 3% of the reaction product is, again, the noxious dioxide which must be further treated by, for instance, a tail gas treating process, in order to reduce the level of SO.sub.2 effluent to within the Environmental Protection Agency standard of not more than 250 parts per million on a dry oxygen-free basis. Consequently, while the Claus process is still a viable alternative for disposing of hydrogen sulfide gas, it has become less attractive because of the cost of carrying out the process and because of the decrease in demand for the elemental sulfur.\nAnother method for treating petroleum to reduce the degrading effects of sulfur is chemical processing to \"sweeten\" sulfur compounds contained in the particular fractions, e.g., the mercaptans which are designated by the formula RSH. \"Sweetening\" denotes that mercaptan sulfur compounds are removed from a refinery stream, or else the mercaptans are converted to less objectionable disulfide compounds, e.g., R-S-S-R, R-S-S-R', etc. A particularly important process employed today is the sweetening of kerosine by the MEROX process to obtain jet fuel.\nWhether sweetening is undertaken by solvent refining processes or by fixed bed adsorption, a caustic solution is generally first used to convert the mercaptan to the ionic state, RS.sup.-. Caustic solution is also helpful in that it removes napththenic acids and other organic acids in general, and other sulfur compounds from cracked petroleum products and petroleum distillate. In fact, caustic treating of petroleum products has been used to improve odor and color nearly as long as the industry itself has been in existence. Numerous equipment modifications and processes have been designed to implement caustic treatment of process streams. Unfortunately, since caustic is quite harmful to organic tissue extreme care must be taken in the use and disposal of alkaline solutions such as sodium hydroxide solutions in order to protect the waterways, rivers, subterranean water formations, and, in many places, the oceans and surrounding seas of industrial areas from caustic pollution. This presents a significant problem to the industry because of the great volume of caustic used in petroleum processing and because all of the solutions used must eventually be discarded as a nonregenerative caustic or as the spent liquor resulting from a regenerative process.\nTo date, the industry generally uses two methods to dispose of spent caustic--neutralization and incineration, incineration being a relatively new trend in waste alkaline liquor disposal. Incineration disposal presents certain advantages over neutralization disposal in that it is, first of all, environmentally cleaner since acid neutralization has a residual OOD (Organic Oxygen Demand) for the naphthenic acids. Other advantages include such benefits as lower capital investment and less operating space required (i.e. ground area). Furthermore, incineration disposal is easier to operate. On the negative side, cost of operation of the incinerator units is high because of the energy required to maintain the elevated temperatures necessary to maintain combustion of the predominantly aqueous solution of alkaline waste. Commercial units presently in operation make use of the combustion of fuel oil, and natural gas to sustain the evaporation of the aqueous parts of the waste liquor and then furnish mostly carbon dioxide, CO.sub.2, to form the innocuous carbonates of sodium or other alkali metals for disposal. Refined fuel oil, and natural gas are very expensive means for disposing of the volumes of caustic discarded as a result of petroleum processing and they are both very valuable commodities in terms of consumer-useable energy sources.\nFrom the foregoing discussion, it can be seen that the elimination of sulfur from petroleum presents several problems including, inter alia, production of alkaline waste liquor and hydrogen sulfide, both of which require either energy and/or expensive treatment to convert them to environmentally safe substances. It is, therefore, a primary purpose of the present invention to provide a process whereby the great volume of alkaline waste liquor which results from petroleum processing can be safely disposed of without expending commercially useable and expensive fuel supplies.\nAnother object of this invention is to provide a process for the formation of innocuous products from potentially harmful and/or noxious sulfur-containing fuels."}
{"text": "Prosthetic implants such as meshes, combination mesh products or other porous prostheses are commonly used to provide a physical barrier between types of tissue or extra strength to a physical defect in soft tissue. However, such devices are often associated with post-surgical complications including post-implant infection, pain, excessive scar tissue formation and shrinkage of the prosthesis or mesh. Excessive scar tissue formation, limited patient mobility, and chronic pain are often attributed to the size, shape, and mass of the implant and a variety of efforts have been undertaken to reduce the amount of scar tissue formation. For example, lighter meshes using smaller fibers, larger weaves, and/or larger pore sizes as well as meshes woven from both non-resorbable and resorbable materials are in use to address these concerns.\nFor treating acute pain and infection, patients with implanted prostheses are typically treated post-operatively with systemic antibiotics and pain medications. Patients will occasionally be given systemic antibiotics prophylactically; however, literature review of clinical trials does not indicate that systemic antibiotics are effective at preventing implant-related infections.\nMany types of soft tissue defects are known. For example, hernias occur when muscles and ligaments tear and allow the protrusion of fat or other tissues through the abdominal wall. Hernias usually occur because of a natural weakness in the abdominal wall or from excessive strain on the abdominal wall, such as the strain from heavy lifting, substantial weight gain, persistent coughing, or difficulty with bowel movements or urination. Eighty percent of all hernias are located near the groin but can also occur below the groin (femoral), through the navel (umbilical), and along a previous incision (incisional or ventral). Almost all hernia repair surgeries are completed with the insertion of a barrier or prosthesis to prevent their reoccurrence. Therefore products used in the management of hernias require some measure of permanent strength. The most commonly employed woven meshes are crafted from polypropylene fibers using various weaves. Tightly woven meshes with the highest strength characteristics and stiffness are very easy for the surgeon to implant; however, there appears to be a positive correlation between the tightness of the weave (correlated to surface area and stiffness), lack of patient mobility, and chronic pain. Newer meshes have larger pore structures and while they are more flexible, they are also more difficult to implant by surgeons. They are extremely difficult for laparoscopic repair, as they have very little recoil associated with them and, when rolled up to insert, they cannot be reflattened and positioned in a quick and efficient manner by the surgeon. Hence, a need still exists for surgical meshes, including hernia meshes, that have sufficient stiffness to facilitate handling and ease of insertion during surgery, yet are or can become sufficiently flexible to be comfortable after implantation.\nSurgical meshes that have been manipulated to improve handling, insertiona and positioning post-insertion are known in the art, but do not employ larger-pore mesh construction. For example, a laparoscopic surgical mesh with extruded monofilament PET coils or rings (e.g., the Bard® Composix® Kugel® hernia patch) increases the overall stiffness of the device and gives a shape memory to the device but does not readily allow for drug loading of the mesh, can not provide temporary stiffening of the mesh component, and can not be further shaped into a fixed three-dimensional structure after manufacture without further processing or alteration. Similarly, meshes with reinforced edges have been produced (e.g., Bard® Visilex®). These meshes have the same disadvantages as those with coils or rings. Additionally, the Kugel patch ring has been reported to break under conditions of use, causing patient morbidity and mortality.\nMeshes produced from a co-weave of a biodegradable material with a non-biodegradable material have been described, e.g., the Johnson & Johnson Vypro® and Vypro II® meshes. In these meshes, polypropylene and polyglactin filaments are braided together before being knitted into a mesh. Such meshes do not change stiffness upon implant as the polyglactin fibers are very fine and flexible. The biodegradable fibers in the VyPro meshes in concert with its particular fiber weave imparts additional flexibility to the mesh such that it distends more easily than the surrounding tissue so that it is more flexible than an equivalent polypropylene fiber mesh with the same weave. Moreover, because the biodegradable polymers of that mesh may be subjected to high temperatures to produce fibers and filaments suitable for weaving, it drastically limits the drugs or biologically active agents that can be included in a biodegradable layer since, under such conditions, the vast majority of biologically-active agents and drugs are unable to withstand the manufacturing temperatures involved in fiber and filament formation. If three-dimensional structures are desired for such meshes, they must undergo further processing to attain such shapes. Finally, these meshes are often more difficult for surgeons to anchor in place because the polyglactin fiber cannot withstand the suturing tension.\nThe present invention overcomes these disadvantages by providing temporarily stiffened and shapeable meshes."}
{"text": "1. Field of the Invention\nThe present invention relates to a memory card having an erroneous erase preventing function and its control apparatus, and a data transmission/reception apparatus.\n2. Description of the Prior Art\nConventionally, an external storage apparatus such as a magnetic tape, magnetic disc, optical disc, memory card and the like is used for storing a data from an electronic apparatus such as computer, digital still camera and the like. Especially, the memory card is widely used for its convenience for carrying and high transfer rate.\nSimilarly as a flexible disc, a memory card as an erroneous erase prevention switch for preventing an erroneous erase of a data stored. However, a conventional erroneous erase prevention switch is provided on a main plane of the card memory body and a user cannot operate the switch when the memory card is mounted on an electronic apparatus.\nIn this case, if the memory card is inserted into an electronic apparatus and the user find that the memory card is disabled for writing, the user cannot write a new data on a data already stored. The user should pull out the memory card from the electronic apparatus, set the erroneous erase prevention switch to a writable mode, and again mount the memory card into the electronic apparatus. It has been desired to improve the switching operationability of the erroneous erase prevention switch."}
{"text": "1. Field of the Invention\nThe present invention relates to truck guard devices and more particularly pertains to a new truck guard device for providing protection from impacts from large animals, other vehicles or other objects to a front end of a semi-truck.\n2. Description of the Prior Art\nThe use of truck guard devices is known in the prior art. U.S. Pat. No. 4,099,760 describes a device for positioning on a vehicle which includes a brush guard and a utility rack. This device allows a vehicle to drive through relatively light brush and prevents damage to the vehicle from the light brush. Another type of truck guard device is U.S. Pat. No. 5,683,128 having a grill assembly that is removably attachable to the front end of a truck for preventing damage to the front end. Yet another such device is found in U.S. Pat. No. 6,290,271 that has a design adapted for positioning on a pick-up truck front end.\nWhile these devices fulfill their respective, particular objectives and requirements, the need remains for a device that is configured for being removably positioned on a semi-truck for preventing damage to the front end of the semi-truck, particularly that damage from impacts with large animals such as deer. Additionally the device should incorporate a tilting means to allow a person to tilt the device away from the front end when a cab of the semi-truck must be lifted during engine maintenance."}
{"text": "The object of the present invention is a method of amplitude-demodulation for radio receivers. Since such receivers are usually executed in digital technology, a digitally realized amplitude demodulator can make use of the timing signals present in the receiver."}
{"text": "The present invention is related to improved medical devices, systems, and methods, with many embodiments being particularly useful for reducing the distance between two points in tissue in a minimally or less invasive manner. Specific reference is made to the treatment of a failing heart, particularly the alleviation of congestive heart failure and other progressive heart diseases. The provided devices, systems, and methods will often be used to resize or alter the geometry of a ventricle in a failing heart, such as by reducing its radius of curvature through the process of excluding a portion of the circumference from contact with blood, and thereby reduce wall stress on the heart and improve the heart's pumping performance. Although specific reference is made to the treatment of congestive heart failure, embodiments of the present invention can also be used in other applications in which tissue geometry is altered.\nExemplary embodiments described herein provide implants and methods for alleviating congestive heart failure and other progressive diseases of the heart. Congestive heart failure may, for example, be treated using one or more implants which are selectively positioned relative to a first wall of the heart (typically an interventricular septum), and another wall of the heart so as to exclude scar tissue and limit a cross sectional area, or distance across a ventricle. Functional deterioration of the heart tissues may be inhibited by decreasing a size of the heart chamber and/or approximating tissues so that stress on the tissues is limited. Implant locations and overall chamber remodeling achieved by placement of a series of implants may be determined so as to provide a beneficial volumetric decrease and chamber shape.\nCongestive heart failure (sometimes referred to as “CHF” or “heart failure”) is a condition in which the heart does not pump enough blood to the body's other organs. Congestive heart failure may in some cases result from narrowing of the arteries that supply blood to the heart muscle, high blood pressure, heart valve dysfunction due to degenerative processes or other causes, cardiomyopathy (a primary disease of the heart muscle itself), congenital heart defects, infections of the heart tissues, and the like. However, in many cases congestive heart failure may be triggered by a heart attack or myocardial infarction. Heart attacks can cause scar tissue that interferes with the heart muscle's healthy function, and that scar tissue can progressively replace more and more of the contractile heart tissue. More specifically, the presence of the scar may lead to a compensatory neuro-hormonal response by the remaining, non-infarcted myocardium leading to progressive dysfunction and worsening failure.\nPeople with heart failure may have difficulty exerting themselves, often becoming short of breath, tired, and the like. As blood flow out of the heart decreases, pressure within the heart increases. Not only does overall body fluid volume increase, but higher intracardiac pressure inhibits blood return to the heart through the vascular system. The increased overall volume and higher intracardiac pressures result in congestion in the tissues. Edema or swelling may occur in the legs and ankles, as well as other parts of the body. Fluid may also collect in the lungs, interfering with breathing (especially when lying down). Congestive heart failure may also be associated with a decrease in the ability of the kidneys to remove sodium and water, and the fluid buildup may be sufficient to cause substantial weight gain. With progression of the disease, this destructive sequence of events can cause the progressive deterioration and eventual failure of the remaining functional heart muscle.\nTreatments for congestive heart failure may involve rest, dietary changes, and modified daily activities. Various drugs may also be used to alleviate detrimental effects of congestive heart failure, such as by dilating expanding blood vessels, improving and/or increasing pumping of the remaining healthy heart tissue, increasing the elimination of waste fluids, and the like.\nSurgical interventions have also been applied for treatment of congestive heart failure. If the heart failure is related to an abnormal heart valve, the valve may be surgically replaced or repaired. Techniques also exist for exclusion of the scar and volume reduction of the ventricle. These techniques may involve (for example) surgical left ventricular reconstruction, ventricular restoration, the Dor procedure, and the like. If the heart becomes sufficiently damaged, even more drastic surgery may be considered. For example, a heart transplant may be the most viable option for some patients. These surgical therapies can be at least partially effective, but typically involve substantial patient risk. While people with mild or moderate congestive heart failure may benefit from these known techniques to alleviate the symptoms and/or slow the progression of the disease, less traumatic, and therefore, less risky therapies which significantly improve the heart function and extend life of congestive heart failure patients has remained a goal.\nIt has been proposed that an insert or implant be used to reduce ventricular volume of patients with congestive heart failure. With congestive heart failure, the left ventricle often dilates or increases in size. This can result in a significant increase in wall tension and stress. With disease progression, the volume within the left ventricle gradually increases and blood flow gradually decreases, with scar tissue often taking up a greater and greater portion of the ventricle wall. By implanting a device which brings opposed walls of the ventricle into contact with one another, a portion of the ventricle may be excluded or closed off. By reducing the overall size of the ventricle, particularly by reducing the portion of the functioning ventricle chamber defined by scar tissue, the heart function may be significantly increased and the effects of disease progression at least temporarily reversed, halted, and/or slowed."}
{"text": "Various heating devices are known to the art of which the following patents are specifically known to the applications, none of which provide the specific features of the present invention: U.S. Pat. Nos. 1,990,640, Doherty; 2,492,865, Huenergardt; 2,687,469, Koci; 2,727,129, Davis, Jr.; 2,740,879, Bieling et al; 3,488,473, Sanders; 3,524,968, Walsh; 3,681,568, Schaefer; 4,096,376, Macklem."}
{"text": "This application is related to U.S. patent application Ser. No. 09/167,392, filed Oct. 6, 1998, entitled xe2x80x9cModular Print Cartridge Receptacle for Use in Inkjet Printing Systemsxe2x80x9d and U.S. patent application Ser. No. 09/167,394, filed Oct. 6, 1998, entitled xe2x80x9cInkjet Printing Systems Using a Modular Print Cartridge Assemblyxe2x80x9d which are herein incorporated by reference.\nThis invention relates to inkjet printers using stationary print cartridges and, more particularly, to servicing stationary print cartridges in an inkjet printing system.\nThermal inkjet hardcopy devices such as printers, graphics plotters, facsimile machines and copiers have gained wide acceptance. These hardcopy devices are described by W. J. Lloyd and H. T. Taub in xe2x80x9cInk Jet Devices,xe2x80x9d Chapter 13 of Output Hardcopy Devices (Ed. R. C. Durbeck and S. Sherr, San Diego: Academic Press, 1988) and U.S. Pat. Ser. Nos. 4,490,728 and 4,313,684. The basics of this technology are further disclosed in various articles in several editions of the Hewlett-Packard Journal [Vol. 36, No. 5 (May 1985), Vol. 39, No. 4 (August 1988), Vol. 39, No. 5 (October 1988), Vol. 43, No. 4 (August 1992), Vol. 43, No. 6 (December 1992) and Vol. 45, No. 1 (February 1994)], incorporated herein by reference. Inkjet hardcopy devices produce high quality print, are compact and portable, and print quickly and quietly because only ink strikes the media.\nAn inkjet printer forms a printed image by printing a pattern of individual dots at particular locations of an array defined for the printing medium. The locations are conveniently visualized as being small dots in a rectilinear array. The locations are sometimes xe2x80x9cdot locationsxe2x80x9d, xe2x80x9cdot positionsxe2x80x9d, or pixelsxe2x80x9d. Thus, the printing operation can be viewed as the filling of a pattern of dot locations with dots of ink.\nInkjet hardcopy devices print dots by ejecting very small drops of ink onto the print medium and typically include a movable carriage that supports one or more print cartridges each having ink ejecting nozzles. The carriage traverses over the surface of the print medium, and the nozzles are controlled to eject drops of ink at appropriate times pursuant to command of a microcomputer or other controller, wherein the timing of the application of the ink drops is intended to correspond to the pattern of pixels of the image being printed.\nThe typical inkjet printhead (i.e., the silicon substrate, structures built on the substrate, and connections to the substrate) uses liquid ink (i.e., dissolved colorants or pigments dispersed in a solvent). It has an array of precisely formed orifices or nozzles attached to a printhead substrate that incorporates an array of ink ejection chambers which receive liquid ink from the ink reservoir. Each chamber is located opposite the nozzle so ink can collect between it and the nozzle. The ejection of ink droplets is typically under the control of a microprocessor, the signals of which are conveyed by electrical traces to the resistor elements. Properly sequencing the operation of each nozzle causes either to eject ink or to refrain from ejecting ink according to the output of the controlling microprocessor to cause characters or images to be printed upon the media as the printhead moves past the media or the media moves past the printhead.\nColor inkjet hardcopy devices commonly employ a plurality of print cartridges, usually two to four, mounted in the printer carriage to produce a full spectrum of colors. In a printer with four cartridges, each print cartridge can contain a different color ink, with the commonly used base colors being cyan, magenta, yellow, and black. In a printer with two cartridges, one cartridge can contain black ink with the other cartridge being a tri-compartment cartridge containing the base color cyan, magenta and yellow inks, or alternatively, two dual-compartment cartridges may be used to contain the four color inks. In addition, two tri-compartment cartridges may be used to contain six base color inks, for example, black, cyan, magenta, yellow, light cyan and light magenta. Further, other combinations can be employed depending on the number of different base color inks to be used.\nThe base colors are produced on the media by depositing a drop of the required color onto a dot location, while secondary or shaded colors are formed by depositing multiple drops of different base color inks onto the same dot location, with the overprinting of two or more base colors producing the secondary colors according to well established optical principles.\nFor many applications, such as personal computer printers and fax machines, the ink reservoir has been incorporated into the pen body such that when the pen runs out of ink, the entire pen, including the printhead, is replaced.\nHowever, for other hardcopy high volume printing applications, such as large format plotting of engineering drawings, color posters and the like, there is a requirement for the use of much larger volumes of ink than can be contained within the replaceable pens. Therefore, various off-board ink reservoir systems have been developed recently which provide an external stationary ink supply connected to the scanning cartridge via a tube. The external ink supply is typically known as an xe2x80x9coff-axis,xe2x80x9d xe2x80x9coff-board,xe2x80x9d or xe2x80x9coff-carriagexe2x80x9d ink supply.\nThere is a trend to use inkjet printing in new specialized printing systems which are very different systems compared to desk-top printers and facsimile machines, or from large format plotters. These specialized printing systems include applications, such as postal printing, postal franking, label printing and bar code printing. Currently, there are no means to design a specialized printing system without a substantial engineering effort.\nIn typical inkjet printers the inkjet print cartridges containing the nozzles are scanned or moved repeatedly across the width of the medium to be printed upon. During this movement across the medium, each of the nozzles is caused either to eject ink or to refrain from ejecting ink according to the program output of the controlling microprocessor. Each completed scan or movement across the medium can print a swath approximately as wide as the number of nozzles arranged in a column of the ink cartridge multiplied times the distance between nozzle centers. After each such completed movement or swath the medium is moved or advanced forward the width of the swath, and the ink cartridge begins the next swath.\nIn inkjet printers the print cartridges need to be periodically serviced. In a scanning carriage printer, a service station is normally located in the scan direction past the edge of the media, since because the scan direction motion is required for printing, it is natural to expand that motion and locate the service station in that direction out of the print zone. Accordingly, when servicing is required, the print cartridges move past the edge of the medium to the location of the service station for servicing.\nIn printing systems which use stationary print cartridges for printing (such as for example, ticket, tag, label and mail printing), there is no scan direction motion because the print cartridges remain in a fixed or stationary position during printing. In addition, in typical stationary printhead printing systems there may be media movement and drive system mechanisms which make accessing the print cartridges difficult. Because of these difficulties, many stationary printhead printing systems do not use service stations. The disadvantage of this approach is that the performance of the printhead decreases as nozzles become dried and ink residue builds up on the printhead orifice plate. Accordingly, periodically it is necessary for an operator technician to manually remove the print cartridges and manually clean the orifice plates. This type of operation is not well controlled and depends on user know-how and consistent execution and exposes the printheads to damage if done incorrectly.\nAccordingly, there is a need for a solution to the servicing of print cartridges in specialty printing systems which often use print cartridges which are stationary during the printing operation with only the media moving through the print zone. With stationary printhead printers, a new means for making the print cartridges accessible to service station components is required.\nThe present invention enables the print cartridges on printers using stationary print cartridges to be accessed by service station components such as wipers, scrapers, cleaning fluid applicators, ink receiving receptacles and cappers. The advantage of the present invention is that enables the print cartridges on stationary printhead printing system to be serviced in a manner similar to what is done in a conventional scanning printhead printing system. The benefits of correctly servicing the inkjet print cartridges are increased printhead quality, increased printhead life, more consistent performance over a wide range of environmental conditions, and reduced operator intervention such as manually clean and replacing print cartridges. The present invention includes a print cartridge support structure for holding one or more print cartridges in a stationary position while applying ink on a media. A printhead is located on each of the print cartridges. The printhead has nozzles for ejecting ink which are arranged in a nozzle array of one or more columns of nozzles. A media movement mechanism provides movement of the media though a print zone located beneath the nozzle array. A printhead servicing station is located outside of the print zone and has one or more servicing modules dedicated for interaction with one of the nozzle arrays when the one of the nozzle arrays is positioned in aligned proximity with its dedicated servicing module. A motorized device coupled to the print cartridge support structure moves the print cartridge support structure out of the print zone to the service station during a period when the print cartridges are not applying ink to the media."}
{"text": "Conventional gold-based catalysts, and particularly nanoporous gold-based catalysts either consist of the nanoporous gold material, or are formed by depositing gold on and/or within metal oxides to form an active structure. While these conventional structures exhibit beneficial activity at relatively low temperatures (e.g. about 0 centigrade), as temperature increases the activity degrades due to the nanoporous gold material forming aggregates (e.g. by sintering), reducing the surface area of the overall structure and thus the catalytic activity thereof.\nThis is particularly undesirable for use in applications to which such catalysts have been earnestly researched. In particular, for an engine (in particular an internal-combustion engine) that operates at a wide range of temperatures, a catalyst that could efficiently treat exhaust gas both in a cold-start phase and at a peak-operation phase when the engine is hot (e.g. several hundred centigrade) would be beneficial."}
{"text": "1. Field of Invention\nThis invention relates to methods and devices for providing a flight crew of an equipped aircraft with timely traffic advice regarding nearby aircraft equipped with operating strobe-type anticollision lights. More specifically, the present invention pertains to methods and devices for detecting anticollision, strobe-like emissions despite background noise and other light pulses which would tend to activate a traffic alert alarm despite the absence of nearby aircraft.\n2. Prior Art\nNumerous devices and methods have been employed over aviation history to assist in the visual detection of nearby aircraft. The use of anticollision strobe lights mounted on the exterior fuselage and wing tips has long been part of standard aircraft equipment for enhancing the ability of flight personnel to see approaching aircraft.\nIt is obviously desirable to facilitate early detection of such anticollision strobe lights so that a pilot may take evasive action well in advance of danger. In daylight hours, however, visual identification of a momentary pulse of light generated at substantial distance is very difficult. Accordingly, scanning devices have been developed to assist in identifying the occurrence of strobe flashes, as well as the approximate azimuth of the strobe light with respect to the aircraft. U. S. Pat. Nos. 3,551,676 and 4,527,158 by Runnels describe the general use of optics and detectors which endeavor not only to identify the strobe pulse of a nearby aircraft, but at the same time provide filtering or blocking of background noise and stray pulses which may generate an alarm signal when in fact aircraft are not in the vicinity. The latter Runnels patent suggests the use of multiple detectors of photosensitive composition as part of a scanning system. Strobe pulse signals are isolated from background noise by use of a bandpass filter which is tuned to the range of frequencies associated with strobe light devices. In addition, a threshold adjustment enables blocking of much of the background noise gathered by the respective detectors.\nU. S. Pat. No. 4,724,312 by Snapper also discloses a detection system embodying forms of filters for eliminating spurious and background signals. In addition, Snapper and Runnels provide a blocking circuit for disabling the detector system during the emittance of pulse radiation from strobes mounted on the monitoring aircraft. These disclosed devices also provide some orientation indication as to the azimuth of the pulse detected. This directional orientation is basically achieved by determining which of a plurality of detectors received the strobe radiation. The azimuth orientation is based on the relative position of that detector on the exterior of the aircraft.\nAlthough the use of passive detectors enhances the ability of a flight crew to detect approaching aircraft having anticollision strobe lights, lack of system efficiency and limited effective range have inhibited commercial acceptance in the small aircraft industry. Such systems may be contrasted with active systems wherein a radar signal or some other form of radiation is emitted by the monitoring aircraft, with subsequent detection of a reflected signal as the basis of detection. In such instances, the detectors are tuned specifically to the frequency of radiation which was emitted, making identification of the reflected signal much less complex.\nIn a passive system such as the present invention, detectors are not tuned to a specific frequency because strobe light frequencies vary, not only in wave length and amplitude, but in pulse rate. Therefore, the system must be prepared to detect a variety of signals which generally conform to a strobe-like flash falling within the frequency range of approximately 300 nm to 1,000 nm with pulse durations of 150 microseconds to 600 microseconds. Pulse repetition rates further vary between 40 and 120 pulses per minute. Accordingly, the challenge of developing an effective passive system for detecting strobe lights is developing hardware and electrical and logic circuitry which can differentiate among these numerous variables for incoming pulse signals.\nA frequent problem with passive strobe detectors has been to develop techniques which isolate true strobe pulses from other signals which have the appearance of a pulse of light similar to that which might be generated by a strobe source. If one designs the range of detection too narrow, true strobe pulses which fall outside the narrow definition may not be identified. For example, hazy or cloudy conditions may affect the detected signal such that it falls outside the narrow range defined as a desired strobe signal. On the other hand, if one broadens the parameters of the system to cover the full range of potential strobe emissions under all weather conditions, the detector system will not only identify such actual strobe pulses, but will trigger false alarms based on reflections of light from ground vehicles and structures, aircraft reflections and changing background such as white clouds against blue sky.\nPast efforts to resolve this dichotomy have generally focused on enhancing the accuracy of detection of a true strobe pulse, to thereby enable distinguishing this pulse from background noise or other similar forms of radiation. Similarly, prior art techniques for enhancing such passive detection systems have concentrated on more accurate definition of the spurious signals which could trigger an alarm condition. The system is than provided with blocking circuits or filters to delete or ignore those specific signals. For example, the Runnels and Snapper approach of disabling the detector during emission of the strobe of the monitoring aircraft is representative of the philosophy of more accurately defining the unwanted signal which may be detected, with subsequent screening of those signals to avoid an erroneous alarm.\nIt has now become apparent to the present inventors that the traditional design strategy of attempting to more accurately specify parameters of (i) the strobe pulse to be detected or (ii) the spurious signal to be screened, filtered or blocked is simply not compatible with the \"real\" world. Too many variables exist in pulse amplitude, width and pulse rate to enable careful refinement of a sufficiently narrow bandpass capable of retrieving only true pulse signals. With respect to spurious signal definition, there are simply too many predictable and unpredictable forms of light pulse to enable development of a catalog of spurious signal forms which the detection system must ignore. Furthermore, some spurious signals may actually have the form and appearance of a true strobe pulse. A system which deletes such a signal form from detection may in fact delete a true pulse in error.\nWhat is needed, therefore, is a passive detection system which is capable of identifying substantially all pulse-type signals similar to an anticollision strobe signal and isolating only those signals which are true strobe pulse signals. Just as important, this system needs to be capable of operating with respect to signals detected at a range of as much as 3 nautical miles, as compared to the more limited ranges of approximately 1 mile for prior art systems."}
{"text": "This invention relates to polyurethane polymers. This invention particularly relates to polyurethane polymers obtained from aqueous polyurethane dispersions.\nPolyurethane dispersions are known, and can be used to obtain polyurethane polymers that can themselves be useful in various applications. Polyurethane/urea dispersion s can be used to obtain, for example: carpet backings; coatings for wood finishing; glass fiber sizing; textiles; adhesives; automotive topcoats and primers; films for packaging; gloves, and other applications. Polyurethane dispersions can be prepared by various processes, including, for example, those described in: U.S. Pat. No. 4,857,565; U.S. Pat. No. 4,742,095; U.S. Pat. No. 4,879,322; U.S. Pat. No. 3,437,624; U.S. Pat. No. 5,037,864; U.S. Pat. No. 5,221,710; U.S. Pat. No. 4,237,264; and, U.S. Pat. No. 4,092,286.\nPolyurethane dispersions can be obtained according to a process that is described in U.S. application Ser. No. 09/039978. Dispersions prepared according to the process described therein can be useful for obtaining polyurethane carpet backings and polyurethane textile backings. Problems can result, however, in the event that a spill occurs wherein a liquid can flow onto, or be absorbed by, the backing on the underside of a carpet. Spilled liquids such as water, urine, beverage drinks, food, blood, and feces can penetrate a carpet backing to the underside of the carpet, which can be inaccessible to various cleaning methods. Polymeric materials having moisture resistance can be useful in protective wear such as gloves, for example, or packaging.\nVarious conventional methods can provide a moisture barrier to polymers obtained from aqueous polymeric dispersions. For example, increased coating weights of a dispersion can be used, thereby increasing the thickness of the polymer. Alternatively, wax added to an aqueous polymeric dispersion can provide a moisture barrier to carpet. Another known method for providing a moisture barrier is to decrease the amount of filler used in a carpet backing formulation. A non-permeable fabric or film can also be applied to a carpet backing, as described in U.S. Pat. No. 5,763,040. Still another method involves applying a fluorochemical on the underside of a secondary backing, and is described in U.S. Pat. No. 5,348,785. Use of fluorochemicals to impart water impermeability is also described in U.S. Pat. Nos. 4,619,853 and 4,643,930. A water-impervious film is described in U.S. Pat. No. 4,336,089. Application of various hydrophobic compositions to a secondary backing is described in U.S. Pat. No. 5,558,916.\nIt would be desirable in the art of preparing polyurethane polymers, to prepare a polyurethane dispersion that includes a component that imparts a moisture barrier to a polyurethane polymer, without adding complexity to a manufacturing process, or significantly changing the other properties of the polyurethane polymer. It would also be desirable in the art to prepare such a polymer by a process that would not significantly increase the cost of manufacture of the polymer, or articles produced therefrom.\nIn one aspect, the present invention is an aqueous polyurethane dispersion comprising an isocyanate terminated prepolymer prepared from a formulation including a polyisocyanate and a polyol mixture, wherein the polyol mixture includes at least one aliphatic mono, di or polyhydroxy compound which has an aliphatic side chain substituent which contains from 5 to 30 carbon atoms and wherein the aliphatic mono, di or polyhydroxy compound which has an aliphatic side chain substituent is present at a concentration of from 0.5 to 10 weight percent of the polymer.\nIn another aspect, the present invention is a moisture resistant polyurethane polymer comprising a polymer layer prepared by applying a layer of an aqueous polyurethane dispersion comprising an isocyanate terminated prepolymer prepared from a formulation including a polyisocyanate and a polyol mixture, wherein the polyol mixture includes at least one aliphatic mono, di or polyhydroxy compound which has an aliphatic side chain substituent which contains from 5 to 30 carbon atoms and wherein the aliphatic mono, di or polyhydroxy compound which has an aliphatic side chain substituent is present at a concentration of from 0.5 to 10 weight percent of the polymer to a substrate and allowing the dispersion to cure.\nIn still another aspect, the present invention is a moisture resistant textile comprising a textile and adherent thereto a polymer prepared by applying a layer of an aqueous polyurethane dispersion comprising an isocyanate terminated prepolymer prepared from a formulation including a polyisocyanate and a polyol mixture, wherein the polyol mixture includes at least one aliphatic mono, di or polyhydroxy compound which has an aliphatic side chain substituent which contains from 5 to 30 carbon atoms and wherein the aliphatic mono, di or polyhydroxy compound which has an aliphatic side chain substituent is present at a concentration of from 0.5 to 10 weight percent of the polymer to the textile and allowing the dispersion to cure.\nMoisture resistant polyurethane of the present invention can be useful in cushioned flooring applications such as attached cushion broadloom, carpet tiles, carpet underlay, or vinyl flooring; adhesives applications; coatings; protective clothing or protective gear such as gloves and aprons; packaging; or any application where moisture resistant polymers can be useful.\nIn one embodiment, the present invention is an aqueous polyurethane dispersion composition that can yield a coating or a foam having good moisture resistance. A polyurethane dispersion of the present invention includes water, and either: a polyurethane; a mixture capable of forming a polyurethane; or a mixture of both. A polyurethane dispersion of the present invention can optionally include: chain extenders; surfactants; fillers; dispersants; foam stabilizers; thickeners; fire retardants, and/or other optional materials that can be useful in a polyurethane formulation.\nThe polyisocyanate component of the formulations of the present invention can be prepared using any organic polyisocyanates, modified polyisocyanates, isocyanate-based prepolymers, and mixtures thereof. These can include aliphatic and cycloaliphatic isocyanates, but aromatic and especially multifunctional aromatic isocyanates such as 2,4- and 2,6-toluenediisocyanate and the corresponding isomeric mixtures; 4,4xe2x80x2-,2,4xe2x80x2- and 2,2xe2x80x2-diphenyl-methanediisocyanate (MDI) and the corresponding isomeric mixtures; mixtures of 4,4xe2x80x2-, 2,4xe2x80x2- and 2,2xe2x80x2- diphenylmethanediisocyanates and polyphenyl polymethylene polyisocyanates PMDI; and mixtures of PMDI and toluene diisocyanates are preferred. Most preferably, the polyisocyanate used to prepare the prepolymer formulation of the present invention is MDI or PMDI.\nA polyurethane dispersion of the present invention is an aqueous dispersion of a polyurethane and/or polyurethane-forming materials. For the purposes of the present invention, polyurethane dispersions can include other polymeric and/or oligomeric compounds including, but not limited to: ureas, biurets, allophonates, and other compounds having functionality that can be derived from the reaction of an isocyanate with the various active hydrogen containing compounds present in the formulation. Polyurethane-forming materials, as the term is used in the present invention, are materials which are capable of forming polyurethane polymers. Polyurethane-forming materials include, for example, polyurethane prepolymers. Prepolymers useful in the practice of the present invention can be prepared by reacting active hydrogen-containing compounds with an amount of isocyanate in excess of the active hydrogen functionality present in the reaction mixture. The isocyanate functionality can be present in an amount of from about 0.2 wt. percent to about 40 wt. percent of the dispersion. A suitable prepolymer can have a molecular weight in the range of from about 100 to about 10,000.\nPrepolymers useful in the practice of the present invention should be substantially liquid under the conditions of dispersal. While it is possible to add a liquid solvent or diluent such as a volatile organic compound (VOC) to the prepolymer to obtain a more fluid prepolymer, it is preferred in this invention to utilize no solvent or diluent, as the removal of these ingredients require extra processing steps and possible environmental consequences from the use of these VOCs. Preferably, the polyurethane dispersions of the present invention are prepared in the substantial absence of an organic solvent.\nActive hydrogen compounds are compounds having functional groups that contain at least one hydrogen atom bonded directly to an electronegative atom such as nitrogen, oxygen or sulfur. Various types of active hydrogen compounds, such as amines, alcohols, polyether polyols, polyester polyols, and mercaptans, for example, are known to those skilled in the art of preparing polyurethane polymers. Active hydrogen compounds suitable for use in the practice of the present invention can be polyols having molecular weights of less than about 10,000.\nThe aliphatic substituted hydroxy functional compound employed in the invention can consist of any compound possessing an aliphatic substituent of from 1 to 30 carbon atoms and one or more hydroxyl groups. Examples of such compounds include glycols such as glycerin, trimethylolpropane, or trimethylolethane which are partially esterified with a 1 to 30 carbon carboxylic acid, preferably a 5 to 22 carbon carboxylic acid, and more preferably a 10 to 15 carbon carboxylic acid. Examples of these compounds include the group consisting of glycerol monooleate, glycerol monobehenate, glycerol monotallate, glycerol monostearate, glycerol monopalmitate, trimethylolpropane monostearate, and mixtures thereof. Other types of compounds include polymeric versions of the aforementioned compounds, and copolymers of these compounds with other glycols, diols and acids or anhydrides. In addition, the aliphatic substituent may be attached to the hydroxyl portion of the compound via heteroatoms such as oxygen or nitrogen, such compounds include the N,N-bis(hydroxyalkyl) alkylamines, one example of which is N-stearyl diethanolamine. Branched chain hydroxy alkanes such as 1,2-dihydroxyoctadecane may also be used. Additionally sorbitan based esters such as the group consisting of sorbitan monostearate, sorbitan monooleate, sorbitan monolaurate, sorbitan monooleate, and mixtures thereof may also be used. Of these, the glycerol monoesters are preferred.\nThe aliphatic substituted hydroxy functional compound employed in the present invention can be used in a suprizingly low concentration to prepare polymers which are effective at repelling moisture. Preferably, the aliphatic substituted hydroxy functional compounds are present at a concentration of from 0.5 to 10 percent of the polymer weight of a polyurethane dispersion. More preferably, the aliphatic substituted hydroxy functional compounds are present at a concentration of from 1.0 to 4.5 percent of the polymer weight of a polyurethane dispersion and even more preferably, the aliphatic substituted hydroxy functional compounds are present at a concentration of from 1.5 to 2.5 percent of the polymer weight of a polyurethane dispersion.\nOther types of aqueous polymeric dispersions can be used in combination with the polyurethane dispersions of the present invention. Suitable dispersions useful for blending with polyurethane dispersions of the present invention include: styrene-butadiene dispersions; styrene-butadiene-vinylidene chloride dispersions; styrene-alkyl acrylate dispersions; or acrylic dispersions; like compounds and mixtures thereof.\nThe present invention optionally includes a chain extender or a crosslinker. A chain extender is used herein to build the molecular weight of the polyurethane prepolymer by reaction of the chain extender with the isocyanate functionality in the polyurethane prepolymer, i.e., chain extend the polyurethane prepolymer. A suitable chain extender or crosslinker is typically a low equivalent weight active hydrogen containing compound having about 2 or more active hydrogen groups per molecule. Chain extenders typically have 2 or more active hydrogen groups while crosslinkers have 3 or more active hydrogen groups. The active hydrogen groups can be hydroxyl, mercaptyl, or amino groups. An amine chain extender can be blocked, encapsulated, or otherwise rendered less reactive. Other materials, particularly water, can function to extend chain length and, therefore, can be chain extenders for purposes of the present invention.\nPolyamines are preferred chain extenders. It is particularly preferred that the chain extender be selected from the group consisting of amine terminated polyethers such as, for example, Jeffamine D-400 from Huntsman Chemical Company, amino ethyl piperazine, 2-methyl piperazine, 1,5-diamino-3-methyl-pentane, isophorone diamine, ethylene diamine, diethylene triamine, aminoethyl ethanolamine, triethylene tetraamine, triethylene pentaamine, ethanol amine, lysine in any of its stereoisomeric forrns and salts thereof, hexane diamine, hydrazine and piperazine. In the practice of the present invention, the chain extender can be used as an aqueous solution.\nWhile optional, in preparing a polyurethane polymer of the present invention use of a chain extender can be advantageous. Generally, a chain extender is employed in an amount sufficient to react with from about zero (0) to about 100 percent of the isocyanate functionality present in the prepolymer, based on one equivalent of isocyanate reacting with one equivalent of chain extender. It can be desirable to allow water to act as a chain extender and react with some or all of the isocyanate functionality present. A catalyst can optionally be used to promote the reaction between a chain extender and an isocyanate.\nCatalysts are optional in the practice of the present invention. Catalysts suitable for use in the present invention include tertiary amines, and organometallic compounds, like compounds and mixtures thereof. For example, suitable catalysts include di-n-butyl tin bis(mercaptoacetic acid isooctyl ester), dimethyltin dilaurate, dibutyltin dilaurate, dibutyltin sulfide, stannous octoate, lead octoate, ferric acetylacetonate, bismuth carboxylates, triethylenediamine, N-methyl morpholine, like compounds and mixtures thereof. An amount of catalyst is advantageously employed such that a relatively rapid cure to a tack-free state can be obtained. If an organometallic catalyst is employed, such a cure can be obtained using from about 0.01 to about 0.5 parts per 100 parts of the polyurethane-forming composition, by weight. If a tertiary amine catalyst is employed, the catalyst preferably provides a suitable cure using from about 0.01 to about 3 parts of tertiary amine catalyst per 100 parts of the polyurethane-forming composition, by weight. Both an amine type catalyst and an organometallic catalyst can be employed in combination.\nThe present invention optionally includes a filler material. The filler material can include conventional fillers such as milled glass, calcium carbonate, aluminum trihydrate, talc, bentonite, antimony trioxide, kaolin, fly ash, or other known fillers. In the practice of the present invention, a suitable filler loading in a polyurethane dispersion can be from about 100 to about 1000 parts of filler per 100 parts of polyurethane. Preferably, filler can be loaded in an amount of at least about 200 pph, more preferably at least about 300 pph, most preferably at least about 400 pph.\nThe present invention optionally includes a filler wetting agent. A filler wetting agent generally can help make the filler and the polyurethane-forming composition compatible with one another. Useful wetting agents include phosphate salts such as sodium hexametaphosphate. A filler wetting agent can be included in a polyurethane-forming composition of the present invention at a concentration of at least about 0.5 parts per 100 parts of filler, by weight.\nThe present invention optionally includes thickeners. Thickeners can be useful in the present invention to increase the viscosity of low viscosity polyurethane dispersions. Thickeners suitable for use in the practice of the present invention can be any thickener known in the art of preparing polyurethane dispersions. For example, suitable thickeners include ALCOGUM VEP-II*(*ALCOGUM(trademark) VEP-II is a trade designation of Alco Chemical Corporation) and PARAGUM 231*(*PARAGUM 231 is a trade designation of Para-Chem Southern, Inc.). Thickeners can be used in any amount necessary to obtain a dispersion of desired viscosity.\nThe present invention can include other optional components. For example, a polyurethane-forming composition of the present invention can include surfactants, blowing agents, frothing agents, fire retardant, pigments, antistatic agents, reinforcing fibers, antioxidants, preservatives, acid scavengers, and the like. Examples of suitable blowing agents include: gases and/or mixtures of gases such as, for example, air, carbon dioxide, nitrogen, argon, helium, and the like; liquids such as, for example, water, volatile halogenated alkanes such as the various chlorfluoromethanes and chlorfluoroethanes; azo-blowing agents such as azobis(formamide).\nFrothing agents are typically introduced by mechanical introduction of a gas into a liquid to form a froth (mechanical frothing). Mechanical frothing of a polyurethane polymer is a procedure known and practice by those skilled in the art of preparing polyurethane polymers. In preparing a frothed polyurethane foam, it is preferred to mix all components and then blend the gas into the mixture, using equipment such as an Oakes or Firestone foamer. In the preparation of a froth for a carpet backing, it is not necessary to obtain a froth that is stable. In a carpet backing production process, a frothed foam typically is spread on the back of a carpet using a spreading tool, which destroys the froth in the process.\nSurfactants are optional, but can be desirable in the practice of the present invention. Surfactants useful herein can be cationic surfactants, anionic surfactants, or a non-ionic surfactants. Examples of anionic surfactants include sulfonates, carboxylates, and phosphates. Examples of cationic surfactants include quaternary amines. Examples of non-ionic surfactants include block copolymers containing ethylene oxide and silicone surfactants. Surfactants useful in the practice of the present invention can be either external surfactants or internal surfactants. External surfactants are surfactants which do not chemically react with the polymer to form a covalent bond during the preparation of the dispersion. Internal surfactants are surfactants which do become chemically reacted into the polymer during dispersion preparation. A surfactant can be included in a formulation of the present invention in an amount ranging from about 0.01 to about 20 parts per 100 parts by weight of polyurethane component. Preferabley, the formualtions of the present invention include polyurethane prepolymers which are not internal surfactants.\nGenerally, any method known to one skilled in the art of preparing polyurethane dispersions can be used in the practice of the present invention to prepare a moisture resistant polymer of the present invention. A moisture resistant polymer of the present invention can be prepared from polyurethane dispersions that are storage-stable or polyurethane dispersions that are not storage-stable. A storage-stable polyurethane dispersion as described herein is any polyurethane dispersion having a mean particle size of less than about 5 microns. A polyurethane dispersion that is not storage-stable can have a mean particle size of greater than 5 microns. For example, a suitable dispersion can be prepared by mixing a polyurethane prepolymer with water and dispersing the prepolymer in the water using a commercial blender. Alternatively, a suitable dispersion can be prepared by feeding a prepolymer into a static mixing device along with water, and dispersing the water and prepolymer in the static mixer. Continuous methods for preparing aqueous dispersions of polyurethane are known and can be used in the practice of the present invention. For example, U.S. Pat. Nos.: 4,857,565; 4,742,095; 4,879,322; 3,437,624; 5,037,864; 5,221,710; 4,237,264; and 4,092,286 all describe continuous processes useful for obtaining polyurethane dispersions. In addition, a polyurethane dispersion having a high internal phase ratio can be prepared by a continuous process as described in U.S. Pat. No. 5,539,021, incorporated herein by reference. In addition, it can be advantageous to combine a continuous process for preparing a prepolymer with a continuous feed dispersion process in order to maximize process efficiency as discussed in pending U.S. application Ser. No. 09/039978.\nThe polymers of the present invention can be applied to any substrate, but preferably the substrate is a textile. More preferably, the substrate is a carpet and the polymer is in the form of a moisture resistant backing. A carpet backing of the present invention can be prepared from a polyurethane dispersion of the present invention, described hereinabove, using either conventional or non-conventional methods in the art of preparing polyurethane-backed carpets. In preparing polyurethane-backed carpets of the present invention, a polyurethane-forming composition can be applied as a layer of preferably uniform thickness onto one surface of a carpet substrate. Polyurethane dispersions of the present invention can be applied as a precoat, laminate coat or as a foam coat. Polyurethane precoats, laminate coats, and foam coats can be prepared by methods known in the art. Precoats, laminate coats and foam coats prepared from dispersions are described in P. L. Fitzgerald, xe2x80x9cIntegral Latex Foam Carpet Cushioningxe2x80x9d, Coat. Fab. 1977, Vol. 7 (pp.107-120), and in R. P. Brentin, xe2x80x9cLatex Coating Systems or Carpet Backingxe2x80x9d, J. Coat. Fab. 1982, Vol. 12 (pp. 82-91), for example.\nA polyurethane-forming composition can be applied to one surface of a carpet ubstrate before it cures to a tack-free state. Alternatively, a polyurethane dispersion containing completely reacted isocyanate functionality can be applied to a suitable substrate, thereby removing the need to cure the polymer. Typically the polyurethane-forming composition is applied to the surface that is attached to a primary backing but can be applied to a secondary backing such as mesh or fleece. The composition can be applied using equipment such as a doctor knife, air knife, or extruder to apply and gauge the layer. Alternatively, the composition may be formed into a layer on a moving belt or other suitable apparatus and dehydrated and/or partially cured, then married to the carpet substrate using equipment such as a double belt (also known as double band) laminator or a moving belt with an applied foam cushion. The amount of polyurethane-forming composition used can vary widely, from about 5 to about 500 ounces per square yard, depending on the characteristics of the textile. After the layer is applied and gauged, water is removed from the dispersion and the layer can be cured using heat from any suitable heat source such as an infrared oven, a convection oven, or heating plates.\nIn the practice of the present invention, any of the steps used in preparing a polyurethane carpet backing can be carried out in a continuous manner. For example, in a first step the prepolymer can be prepared from a suitable active hydrogen containing compound in a continuous manner; the prepolymer can be fed, as it is obtained in the first step, into a mixing device with water to obtain an aqueous dispersion; the aqueous dispersion can be applied to a carpet substrate in a continuous manner to obtain a polyurethane backed carpet."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and an apparatus for producing a semiconductor device which is capable of obtaining a crystalline silicon film high in uniformity by improving the flatness of a glass substrate during producing an insulated gate semiconductor device such as a thin film transistor (TFT) formed using a non-single crystal silicon film which is formed on the substrate, or other semiconductor devices. In particularly, the present invention is useful in producing a semiconductor device formed on the glass substrate.\n2. Description of the Related Art\nIn recent years, several researches have been made into insulated gate field effect transistors having a thin film-shaped active layer (so-called active region) on an insulated substrate, which are so-called TFT).\nThose TFTs are classified, for example, into an amorphous silicon TFT and a crystalline silicon TFT, depending upon the material or the crystal state of a semiconductor as used. The crystalline silicon stated here is directed to non-single crystal silicon which is not single crystal. Hence, those TFTs are generally called “non-single crystal silicon TFTs”.\nIn general, the electric field mobility of a semiconductor which is in an amorphous state is small, and therefore not available to the TFTs that require a high speed operation. Also, the amorphous silicon cannot be used to produce a p-channel TFT (PMOS TFT) since the p-type electric field mobility of the amorphous silicon is remarkably small, and thus a complementary MOS circuit (CMOS) cannot be formed by the combination of the p-channel TFT and an n-channel TFT (NMOS TFT).\nThe crystalline semiconductor is larger in the electric field mobility than the amorphous semiconductor, and therefore enables a high speed operation. The crystalline silicon can be used for obtaining not only the NMOS TFT but also the PMOS TFT, thereby forming a CMOS circuit.\nA crystalline silicon film is obtained by thermally cooling an amorphous silicon film obtained through the vapor phase growth technique at an appropriate temperature (600° C. or higher) for a long period of time, or by irradiating an intense light such as a laser beam (optically annealing).\nHowever, in the event of using a glass substrate which is inexpensive and rich in processability as an insulating substrate, it is extremely difficult to obtain, by only annealing, a crystalline silicon film which is satisfactorily high in electric field mobility (high to the degree that the CMOS circuit can be formed).\nThis is because the above glass substrate is generally low in strain point (about 600° C.), and thus when the temperature of the substrate is elevated up to a temperature required for obtaining a crystalline silicon film sufficiently high in mobility, the substrate is warped.\nIn the event of applying the optically annealing technique for crystallizing a silicon film formed on the glass substrate, a high energy can be applied only to the silicon film without elevating the temperature of the substrate so much. Hence, the optically annealing technique is very effective in crystallizing the silicon film formed on the glass substrate.\nIt has been found that a high power pulse laser such as an excimer laser is the most suitable for the optimum light source for optically annealing. The maximum energy of the laser is extremely larger than that of the continuously oscillating laser such as an argon ion laser, and the mass productivity could be enhanced using a large spot which is several cm2 or more in size.\nHowever, because the beam as normally used is square or rectangular in shape, the beam is required to move vertically and horizontally for processing a single substrate having a large area. Thus, the optically annealing technique needed to be still improved from the viewpoint of the productivity.\nThe above matter could be remarkably improved with a technique in which the beam is deformed linearly, the width of the beam is a length that exceeds that of a substrate to be processed, and the substrate is scanned relatively by the beam (The scanning operation in the specification means that linear laser beams are irradiated onto the substrate while being overlapped one on another bit by bit.). The details are disclosed in Japanese Patent Unexamined Publication No. 5-112355.\nThe silicon film which is still high in crystallinity can be prepared by conducting the thermally annealing process prior to conducting the optically annealing process. In the thermally annealing process as described in Japanese Patent Unexamined Publication No. 6-244204, utilizing the effect that an element such as nickel, iron, cobalt, platinum, paradium (hereinafter referred to as “crystallized catalytic element, or simply “catalytic element”) promotes the crystallization of amorphous silicon, the crystalline silicon film can be obtained by the thermally annealing process at a lower temperature for a shorter period of time in comparison with the normal case.\nTFTs arranged in the form of a matrix are formed with a crystalline silicon film formed using the thermally annealing process and the optically annealing process together, and the distribution of their threshold voltage in the substrate surface is investigated.\nFIG. 2 shows the distribution of the threshold values of the TFT using the crystalline silicon film formed through the conventional method, within the substrate surface. The distribution is U-shaped as shown in FIG. 2.\nFIG. 4 shows the arrangement of TFTs on the glass substrate. The data in FIG. 2 is obtained, as shown in FIG. 4, in such a manner that the TFTs of 400×300 pieces are arranged in the form of a matrix in a region of 40×50 mm on a Corning 1737 substrate of 100 mm2, and the respective locations of 400 TFTs disposed laterally in a line from one end to the other end (a portion surrounded by a dotted line in FIG. 4) are indicated correspondingly in the axis of abscissa.\nWhen the pixel matrix that constitutes the pixel portion of a liquid crystal display has the distribution of threshold voltages shown in FIG. 2, the display state becomes nonuniform, resulting in a defective image.\nAs a result of researching the cause that the threshold voltage exhibits such a U-shaped distribution within the substrate surface by the applicant, it has been found that a tendency of the U-shaped distribution is very similar to the warp of the substrate immediately before a laser beam is irradiated onto the substrate.\nAlso, no warp of the substrate is found in the glass substrate immediately after an amorphous silicon film is formed on the glass substrate, and it has been found that the warp of the substrate is caused because, during a heat treatment (by which the film grows in the solid phase into a crystallized film) subsequent to the amorphous silicon film forming process, a silicon film (or silicon oxide film) is contracted higher than that of the glass substrate in cooling the substrate after the heat treatment. The warp of the substrate is produced in a U-shape viewed from the film formed on the substrate.\nFIG. 3 shows a state in which a laser annealing is conducted on the silicon film formed on the substrate which has been warped. From FIG. 3, when the laser annealing is conducted on the silicon film in such a state where the substrate is warped, the focal point of the laser beam is shifted differently at the respective locations on the substrate.\nIt is presumed that the shift of the focal point makes the crystallinity of the silicon film different from each other within the substrate surface, so that the threshold voltage exhibits a specified distribution within the substrate surface.\nThe warp of the substrate immediately before a laser beam is irradiated onto the substrate having 100 mm square is different by about 50 μm between the central portion and the edge portion of the substrate. The degree of the warp fell within a range of about 20 to 200 μm although it depends upon the temperature of the above heat treatment process, a time necessary for processing, the material of the substrate, or the like. There is a case in which when the size of the substrate is about 500 mm square, its warp becomes about 1 to 2 mm."}
{"text": "1. Field of the Invention\nThe present invention relates to a resistive memory device having a storage element whose data storage state changes based on resistance change dependent on the application voltage, and an operating method thereof.\n2. Description of the Related Art\nThe resistive memory device is being actively used in resent years because it does not demand the electrical power for holding stored data. The floating gate (FG) flash memory, which is mainly employed as the memory device currently, has a disadvantage of low data writing speed. In the FG flash memory, not only that the data writing speed itself is low, but also that verify operation is needed in data writing. Therefore, the total writing time is long.\nFor resolution of this disadvantage, there are methods in which the verify operation is improved to shorten the total writing time. As one of the methods, there has been proposed a technique in which so-called direct verify operation is applied to a resistive memory device, which has higher writing speed (refer to e.g. Japanese Patent Laid-open No. 2007-133930 (hereinafter, Patent Document 1)).\nIn normal verify operation, for a memory cell, the bit line is pre-charged after the end of a recording process (program), and thereafter verify readout (detection) operation is carried out. A series of operation of repeating the verify operation composed of the program and the detection until the desired resistance change is obtained is carried out. This series of operation will be referred to as the “writing verify operation” if it is for writing, and will be referred to as the “erasing verify operation” if it is for erasing. Hereinafter, the term “recording process” or “program” will refer to recording (program) of data of “1” or “0” no matter whether it is for writing (changing any storage logic from the initial state) or for erasing (returning the storage logic to the initial state).\nIn contrast, in the direct verify operation, after the end of a recording process, verify readout (detection) is performed by directly using the charge left in the bit line at this time to thereby shorten the total time. In the direct verify operation, the total writing operation time can be shortened corresponding to the absence of pre-charge for verify readout.\nThe writing verify operation described in Patent Document 1 will be described below.\nFIG. 1 shows how the bit line voltage (hereinafter, referred to as the BL voltage) changes in a recording (writing) process for writing and the subsequent readout (detection) process when the writing verify operation is carried out by the method described in Patent Document 1.\nThe potential change of the bit line shown in FIG. 1 will be described below.\nBy the writing operation, the BL voltage changes from a common potential Vcommon to a writing voltage Vwrite. When the writing by this writing operation has succeeded, the resistance of the variable resistance element has become low. Thus, when the supply of the writing voltage to the bit line is stopped, the BL voltage changes to a certain sufficiently-low potential corresponding to the resistance of the variable resistance element.\nOn the other hand, when the writing has failed, the resistance of the variable resistance element has not become sufficiently-low but is high. Thus, the BL voltage does not much decrease although the supply of the writing voltage is stopped.\nIf a reference potential Vref-write for writing is set between these potentials in advance, the output dependent on whether the writing has resulted in success or failure is obtained by a sense amplifier, and thus whether the writing has been correctly performed can be checked.\nFIG. 2 shows how the BL voltage changes in a recording (erasing) process for erasing and the subsequent readout (detection) process when the erasing verify operation is carried out by the method described in Patent Document 1.\nIn the method of the related art, the so-called pre-charge process, in which voltage is applied to the bit line BL in order to carry out the verify operation, is carried out in the erasing process, in which an erasing voltage Verase is applied to the bit line BL for erasing.\nThe potential change of the bit line BL dependent on the state as to the success or failure of erasing is different from the potential change in the writing verify operation.\nSpecifically, after the BL voltage is lowered from the common potential Vcommon to the erasing voltage Verase in the erasing process, if the erasing has succeeded, the resistance of the variable resistance element has changed to high resistance, and thus the BL voltage does not much change although the erasing signal is turned off.\nIf the erasing has failed, the resistance of the variable resistance element has not become sufficiently-high but is low, and therefore the BL voltage rises up to a certain potential.\nIf a reference potential Vref-erase for erasing is set in advance between these two potentials taken as the value of the BL voltage, the output dependent on the success or failure of the erasing is obtained by the sense amplifier. Thus, whether the erasing has been correctly performed can be checked.\nBy using the process in which voltage is applied to the bit line BL to perform writing or erasing for the variable resistance element in the selected cell also as the pre-charge process for the verify operation in this manner, the time demanded for the writing verify operation and erasing operation is shortened."}
{"text": "The present invention relates to a titanium dioxide pigment having excellent light fastness and excellent hydrophobic and dispersible properties, and a method for production thereof. More particularly, the invention relates to a titanium dioxide pigment suited for use in plastics and a method for production thereof.\nTitanium dioxide pigments are generally poor in light fastness. When titanium dioxide pigments are used in plastics, paints, ink, etc, these pigments thus promote color change, fading or deterioration of organic components such as resin, oil and fats, etc., upon exposure to UV rays. For this reason, titanium dioxide pigments are normally coated with hydrated compounds of aluminum, silicon, zirconium, etc. on the surface of pigment particles thereby to improve light fastness.\nFurther since titanium dioxide pigments have a hydrophilic nature and when these pigments are incorporated into plastics followed by processing especially into a thin layered film of polyolefin type such as polyethylene or polypropylene, the phenomenon of void formation, which is a so-called lacing problem ascribable to moisture contained in titanium dioxide pigments, occurs to cause deterioration in processing properties such as dispersibility or high packing.\nJP, 60-3430, B discloses titanium dioxide pigments having excellent hydrophobicity and dispersibility, which are coated with organosilicons such as silane coupling agents, polysiloxanes, etc., polyhydric alcohols and hydrated oxides of aluminum.\nAccording to this technique, titanium dioxide pigments are imparted hydrophobicity to improve processability, but light fastness is yet insufficient. In order to impart light fastness to titanium dioxide pigments, it is necessary to coat the surface of titanium oxide particles with a sufficient amount of the hydrated oxides of aluminum, silicon, zirconium, etc. Then, the moisture content in titanium dioxide pigments increases to reduce processability. That is, known titanium dioxide pigments could not satisfy light fastness and processing properties (lacing resistance, etc.) at the same time. Furthermore, in order to remove bound water from these hydrated oxides and reduce the moisture content in the resulting anhydrides, heating and firing at high temperatures are required. However, heating and firing at high temperatures result in problems that dispersibility decreases and production costs increase.\nIn order to overcome the problems of the prior art described above, the present invention has been made to provide titanium dioxide pigments having excellent light fastness, minimized change in color, highly hydrophobic property and excellent processability including dispersibility, high packing, etc., with less lacing, and to provide a method for producing the same.\nThe present inventors have made extensive studies to solve these problems. As a result, the inventors have found that by using an aluminum phosphate compound in place of the hydrated oxides of aluminum, silicon or zirconium described above, light fastness of titanium dioxide pigments can be improved even by coating with an aluminum phosphate compound in a small amount, and further by coating with the hydrolyzate of an organosilane compound, titanium dioxide pigments can be rendered hydrophobic to reduce the water content in the resulting titanium dioxide pigments so that titanium dioxide pigments having excellent light fastness and processability for plastics can be obtained. Based on the finding, the present invention has been accomplished.\nThat is, the present invention is directed to a titanium dioxide pigment comprising titanium dioxide particles, a coating layer containing an aluminum phosphate compound and a coating layer containing the hydrolyzate of an organosilane compound, formed on the surface of the particles, and characterized in that a difference in water content between 100xc2x0 C. and 300xc2x0 C. as determined by the Karl Fischer method, is not greater than 1500 ppm. The present invention is also directed to a method for producing the titanium dioxide pigments.\nThe titanium dioxide particle that is the substrate of the titanium dioxide pigment of the present invention preferably has an average particle diameter of 0.1 to 0.4 xcexcm determined by electron microscopic photography. The crystal form of the titanium dioxide particle may be anatase structure, rutile structure or mixture of anatase structure and rutile structure. Any method is available for preparing the titanium dioxide particles without any limitation. The titanium dioxide particles may be those obtained by a so-called sulfate process which comprises hydrolyzing a titanium sulfate solution, or by a so-called chloride process which comprises gas phase oxidation of a titanium halide.\nThe titanium dioxide pigment of the present invention has acquired excellent light fastness by providing the coating layer containing the aluminum phosphate compound.\nThe aluminum phosphate compound may be any one of aluminum orthophosphate and aluminum metaphosphate. The coating amount of the aluminum phosphate compound ranges preferably from 0.1 to 2.0 wt % in terms of AlPO4 based on titanium dioxide as its substrate. With the amount smaller than the lower limit, the intended light fastness is not obtained and when the amount exceeds the upper limit, the bound water contained in the aluminum phosphate compound causes foaming in plastic molding such as processing of thin layer films, etc., which results in deterioration in lacing resistance.\nThe titanium dioxide pigment of the present invention has acquired excellent processability, dispersibility and hydrophobicity by further providing the coating layer containing the hydrolyzate of organosilane compound.\nThe organosilane compound is desirably represented by general formula (1):\nRnxe2x80x94Sixe2x80x94(ORxe2x80x2)4xe2x88x92nxe2x80x83xe2x80x83(1)\nwherein R is a hydrocarbon group having carbon atoms of not greater than 10 containing at least one of an alkyl, vinyl and methacryl; Rxe2x80x2 is methyl or ethyl; and n is an integer of 1 to 3, provided that when n is 2 or 3, R may be the same or different hydrocarbon group. When the carbon atom number of the hydrocarbon group R is 11 or larger, not only it is difficult to cause hydrolysis but heat resistance of such a hydrolyzate-coated titanium dioxide pigment becomes poor, titanium dioxide powders become yellow upon drying after the coating treatment and heating during grinding, which leads to yellowing of plastics molded therefrom.\nThe hydrolyzate referred to in the present invention is used to mean the silanol products obtained via hydrolysis of the alkoxy group in the organosilane compounds shown by the formula above, the oligomers or polymers containing siloxane bonds by further polycondensation between the silanol compounds, and a mixture of at least two compounds selected from these compounds. The hydrolyzate may contain unreacted silane compounds. In view of efficient hydrophobicity, it is particularly preferred that the coating be finally in the form of a polymer. When the coating treatment is performed in aqueous slurry as described below, however, it is preferred to use a hydrophilic silanol monomer as a treating agent added to an aqueous slurry, the hydrophilic silanol monomer is coated on the titanium dioxide then followed by polycondensation.\nSpecific examples of the organosilane compound described above are methyltriethoxysilane, trimethylethoxysilane, methyltrimethoxysilane, vinyltriethoxysilane, vinylmethyldiethoxysilane, vinyltrimethoxysilane, n-butyltrimethoxysilane, n-butylmethyldimethoxysilane, n-butyltrimethoxysilane, isobutyltriethoxysilane, isobutyltrimethoxysilane, hexyltriethoxysilane, hexylmethyldimethaxysilane, hexyltrimethoxysilane, octyltriethoxysilane, octyltrimethoxysilane, desyltriethoxysilane, desyltriethyoxysilane, methacryloxypropyltriethoxysilane, methacryloxypropyltrimethoxysilane, methacryloxypropylmethyldimethoxysilane, etc.\nThe amount of the hydrolyzate of the organosilane compound is preferably 0.05 to 3.0 wt %, more preferably 0.2 to 2.0 wt %, when calculated as the organosilane compound, based on the titanium dioxide substrate. When the amount is less than the lower limit, excellent hydrophobicity, dispersibility or lacing resistance are not obtained since the blocking effect of hydroxyl group on the surface of titanium dioxide particles is poor. On the other hand, even with the amount larger than the range above, the effect corresponding to the addition amount of the organosilane compound is not obtained, and such is disadvantageous in an economical aspect.\nEither the coating layer containing aluminum phosphate or the coating layer containing the hydrolyzate of the organosilane compound may be formed on the side closer to the titanium dioxide particles. Preferably, the surfaces of titanium dioxide particles are entirely covered by these layers. However, the titanium dioxide pigment of the present invention may contain uncoated titanium dioxide particles, partly coated titanium dioxide particles or titanium dioxide particles covered either by the coating layer containing aluminum phosphate or by the coating layer containing the hydrolyzate of the organosilane compound, within such a range that does not damage the purpose of the invention.\nPreferably, the titanium dioxide pigment of the present invention comprises titanium dioxide particles, the coating layer containing aluminum phosphate and the coating layer containing the hydrolyzate of the organosilane compound. However, the titanium dioxide pigment may have a coating layer other than the coating layer containing aluminum phosphate and the coating layer containing the organosilane compound hydrolyzate, e.g., a coating layer containing the hydrated compounds of aluminum, silicon, zirconium, etc. hitherto known within such a range that does not damage the purpose of the invention. Location of these layers is not particularly restricted, either.\nThe titanium dioxide pigment of the present invention should be that a difference in water content between 100xc2x0 C. and 300xc2x0 C., determined by the Karl Fischer method described below, is not greater than 1500 ppm (water content by the Karl Fischer method). When the difference in water content by the Karl Fischer method exceeds the range above, lacing resistance or high packing property becomes poor.\nMethod for Determining the Water Content by the Karl Fischer Method:\nAfter titanium dioxide was equilibrated by allowing to stand at constant temperature of 25xc2x0 C. and constant humidity of 55% relative humidity, 1 g of the sample was determined in terms of the water content by the Karl Fischer method at 100xc2x0 C. and 300xc2x0 C., using a Karl Fisher device for determining the water content and a water vaporizer attached to the device (both manufactured by Mitsubishi Chemical Corporation) to calculate the difference.\nIn the present invention, it is preferred to set the dispersibility at 20 kg/cm2 or less. The dispersibility referred to in the present invention is expressed by an increase in resin pressure determined by the following method.\nMethod for Evaluation of Dispersibility:\nTitanium dioxide pigment, 500 g, is blended with 500 g of frozen and ground polyethylene resin (Sumikasen L-750, made by Sumitomo Chemical Co., Ltd.) and 20 g of zinc stearate with a juice mixer for 5 minutes. Using a Labo Plastomill twin screw extruder which is made by Toyo Seiki Seisaku-sho, Ltd. and set to have a resin temperature of 280xc2x0 C. and a screen of 1450 mesh in the discharge side, the blend is melt-extruded over an hour. Resin pressures are measured at the start of extrusion and an hour after the extrusion, and its difference is made an increased resin pressure.\nFurthermore, the present invention is directed to the method for producing the titanium dioxide pigment which comprises the step of preparing an aqueous titanium dioxide slurry containing titanium dioxide particles, the step of forming the coating layer containing the aluminum phosphate compound and the step of forming the coating layer containing the hydrolyzate of the organosilane compound, on the surface of titanium dioxide particles. The steps of forming the coating layer containing the aluminum phosphate compound and forming the coating layer containing the hydrolyzate of the organosilane compound, on the surface of titanium dioxide particles, may be in a reversed order.\nThe solid content of titanium dioxide in the aqueous slurry is in the range of preferably 50 to 800 g/liter, more preferably 100 to 500 g/liter. With the content higher than the upper limit, the viscosity of the aqueous slurry increases overly to make a uniform coating on the surface of titanium dioxide particles difficult. With the content lower than the range above, operability on an industrial scale becomes poor.\nThe coating layer containing aluminum phosphate is formed by adding an aqueous solution containing phosphoric acid and aluminum compounds into the titanium dioxide slurry and adjusting pH to 4.5 to 8.0, preferably 4.5 to 7.0, with an aqueous solution of a basic compound. When the pH is higher than 8.0, aluminum oxides are formed in a different phase in which the phosphoric acid compound is not present; in this case, the coating layer of the aluminum phosphate compound contemplated in the present invention is not obtained since the phosphoric acid is not stabilized. When the pH is below 4.5, aluminum phosphate, even if once formed, will be dissolved again. It is advantageous for the coating layer formation to maintain the slurry at a pH of 0 to 4.0 until the pH is adjusted to the range described above. The aluminum phosphate layer can be placed either as a first layer on the surface of titanium dioxide or as a second layer after coating of the hydrolyzate of the organosilane compound, which will be later described. The titanium. dioxide pigment may also be treated in the order of first adding either one of a phosphoric acid compound aqueous solution and an aluminum compound aqueous solution and then adding the organosilane compound to form the coating layer of the hydrolyzate of the organosilane compound, then adding the other, adjusting pH and finally forming the aluminum phosphate coating thereon. Also, the pigment may be treated by previously adding either the phosphoric acid compound aqueous solution or the aluminum compound aqueous solution, next adding the organosilane compound, then adding the other aqueous solution, adjusting pH to form the aluminum phosphate coating layer, and then forming the coating layer of the hydrolyzate of the organosilane compound.\nAny phosphoric acid compounds and aluminum compounds may be usable without any particular restriction, so long as these compounds are soluble in water. Examples of the phosphoric acid compound include orthophosphoric acid, metaphosphoric acid, pyrophosphoric acid and salts thereof. Examples of the aluminum compound are sodium aluminate, aluminum sulfate, aluminum nitrate, and the like. For the pH adjustment, there may be employed acidic compounds, e.g., inorganic acids such as sulfuric acid, hydrochloric acid, etc., or organic acids such as acetic acid, formic acid, etc., and inorganic basic compounds such as sodium hydroxide, potassium hydroxide, ammonia, etc.\nPreferred amounts of the phosphoric acid compounds and aluminum compounds used are appropriately chosen to form 0.1 to 2.0 wt % of aluminum phosphate calculated as AlPO4, based on titanium dioxide as its basis.\nNext, the hydrolyzate of the organosilane compound can be formed, e.g., by any one of (1) a method which comprises adding the organosilane compound to a titanium dioxide aqueous slurry having an acidic pH value adjusted to 0.5 to 6, preferably 2 to 4; and (2) a method which comprises adjusting pH of the aqueous dispersion of the organosilane compound previously prepared to an acidic range of 0.5 to 6, preferably 2 to 4. When the pH is in the range described above, hydrolysis of the organosilane compound proceeds, whereby hydrophilic silanols tend to be formed and hydrophobic polycondensates of silanols are formed only with difficulty. Thus, handling in an aqueous slurry becomes easy and such is preferred. Alternatively, the hydrolyzate may also be formed by preparing silanol monomers, coating the monomers on the titanium dioxide particles, adjusting pH of the slurry to a neutral range or higher, and performing polycondensation to form oligomers or polymers on the surface of the titanium dioxide particles.\nJP, 5-221640, A and WO 95/23194 disclose a method which comprises adding an organosilane compound to an alkaline aqueous dispersion of titanium dioxide, which is different from the method of the present invention. However, when the aqueous slurry of titanium dioxide or the aqueous dispersion of an organosilane compound is inclined toward a neutral or alkaline pH region, a rate of the hydrolysis decreases or polycondensation occurs between silanols to make processing in an aqueous slurry difficult. In addition, the reactivity with titanium dioxide is lowered, which leads to failure of sufficient adherence to the surface of titanium dioxide particles to cause such a problem that the desired efficiency is not obtained.\nThe amount of the organosilane compound employed is appropriately chosen to produce the hydrolyzate in an amount of 0.05 to 3.0 wt %, preferably 0.2 to 2.0 wt %, based on titanium dioxide.\nThe temperature for the hydrolysis is not particularly limited but preferably ranges from 5 to 40xc2x0 C. When the temperature is higher than the range above, the organosilane compound, which does not undergo the hydrolysis, tends to readily evaporate or volatilize so that a fixing rate of the hydrolyzate of the organosilane compound onto the surface of titanium dioxide particles becomes low.\nThe coating of the titanium dioxide particle surface with the hydrolyzate of the organosilane compound can be effected, e.g., by elevating the temperature for the hydrolysis to that for the coating treatment. The temperature for the coating treatment is preferably between 45 and 100xc2x0 C., more preferably between 60 and 100xc2x0 C. When the temperature is lower than the range above, it is difficult to perform the coating, and the temperature higher than the range above makes industrial operations under normal pressure difficult.\nIt is desirable to prepare a dispersion by previously dilute the organosilane compound with a mixture of water and a lower alcohol, or to prepare a slurry of titanium dioxide using a mixed solution of water and a lower alcohol, because the dispersibility of the organosilane compound in the dispersion or slurry is improved.\nExamples of the lower alcohol include methyl alcohol, ethyl alcohol, propyl alcohol, isopropyl alcohol, butyl alcohol, etc. The mixing ratio of the lower alcohol can be suitably set forth.\nAfter the steps of coating the aluminum phosphate compound and the hydrolyzate of the organosilane compound are completed, the slurry is filtered, fractionated, dried and subjected to dry grinding to recover the titanium dioxide pigment. By the heat treatment upon drying, polycondensation between the hydrolyzates of the organosilane compound is promoted to reinforce the hydrophobicity. Therefore, the drying temperature is set preferably at 50 to 250xc2x0 C., more preferably at 100 to 200xc2x0 C. When the heat treatment is carried out at the temperature higher than the range above, the coated organosilane compound causes an undesirable change in color. Any dry grinding is available from conventional dry grinding including air pulverization, impact pulverization, etc., and has no particular restriction.\nThe present invention is further directed to a resin composition which comprises formulating the titanium dioxide pigment described above, which is almost free from lacing, pinholes, etc. or poor processability that poorly dispersed particles of the titanium dioxide pigment are projected out of the surface. The resin composition has excellent surface smoothness and glossy surface, and excellent light fastness almost free from change, fading and deterioration in color, etc. under UV rays. A kneader generally available can be used to formulate the titanium dioxide pigment in plastic resin. Examples of the kneader include a twin screw extruder, a Banbury mixer, etc. As plastic resins used in the resin composition of the present invention, there are thermosplastic resins such as polyolefin resin, vinyl chloride resin, polystyrene resin, ABS resin, engineering plastics, etc., thermosetting resins such as phenol resin, urethane resin, unsaturated polyester resin, etc.; various resins can be used without any particular limitation. A proportion of the titanium dioxide pigment in plastic resin is not particularly restricted but the titanium dioxide pigment is preferably formulated in an amount of 0.01 to 900 parts by weight, more preferably 0.1 to 200 parts by weight, based on 100 parts by weight of plastic resin. A variety of additives known to one skilled in the art as additives to resin compositions may be incorporated depending upon use."}
{"text": "The present invention pertains to thrombus filters. In particular, the present invention relates to thrombus filters which can be securely affixed at a selected location in the vascular system and removed when no longer required.\nIn recent years, pulmonary embolism has become an increasingly common medical emergency. Pulmonary embolisms may be caused by venous thrombosis, which in turn may be caused by blood flow retention, venous intima damage or coagulation abnormalities. Emergency or prophylactic treatments for these conditions includes anti-coagulant therapy, thrombolytic therapy, thrombectomy and inferior vena cava blocking procedures.\nAmong these therapeutic options, when an inferior vena cava blocking procedure is selected, one option is to perform a laparotomy under general anesthesia during which the inferior vena cava is ligated, sutured and shortened, or clipped. A laparotomy however, requires a general anesthetic and is susceptible to thrombosis formation due to the discontinuation of anti-coagulate therapy prior to surgery. Another option is to intravenously insert a thrombus filter into the inferior vena cava which requires a local anesthetic. Percutaneous filter insertion is now widely employed since it requires only a local anesthetic. However, such filters become affixed in the inner vena cava wall by neointimal hyperplasia within two or three weeks after being implanted, after which they cannot be removed percutaneously.\nPrior thrombus filters which may be removed percutaneously have included generally linear or column-like struts which tend to transfer forces along the longitudinal axis of the struts. If the vessel is compressed asymmetrically, the filter elongates forcing the struts into the vessel wall which may damage the vessel."}
{"text": "1. Field of Invention\nThe present invention relates to self-contained inertial navigation systems (INS) for interactive control using movable controllers in applications like computer display games.\n2. Related Art\nThe Nintendo Wii Remote™ wireless controller is an example of the most recent state of the art advances in user interactive controllers for computer display game systems. It is a movable wireless remote controller, which is hand-held by the interactive user, that transmits input data to the computer controlled game display system via conventional short range wireless RF transmissions e.g., a Bluetooth™ system, and receives data via infra-red light sensors. This game controller is described in detail in Published U.S. Application US2007/0060384, (Mar. 15, 2007).\nWith movable controllers for game systems like the Nintendo Wii, it is desirable to use a self-contained INS system to sense and track the relative linear and angular motion of the movable controller. Current state of the art movable controllers do not have this ability. For example, the Wii Remote can use data it receives via its infra-red light sensors to infer information about its position and orientation from a set of external infra-red light sources that have been placed in the environment in some known configuration. But the use of light sensors means that the device depends on the light sources and is not, therefore, self-contained. The use of external signal sources is burdensome because the user must set up and configure those external sources. Furthermore, the user must restrict movements made with the controller so as to keep those external sources in view. A self-contained system has no such restriction on movement and requires no setup or configuration of external sources by the user.\nSelf-contained INS systems typically use sensors like accelerometers and gyroscopes. State of the art movable controllers like the Wii Remote use a tri-axial accelerometer. However, a single tri-axial accelerometer is insufficient to calculate all six degrees of freedom required to infer the linear and angular motion of a movable controller. In fact, it is impossible to even determine whether the controller is being translated or rotated, since a fixed rotation and a linear acceleration could generate the same set of readings on a single tri-axial accelerometer. However, by making assumptions on how the controller is held and along which axis it will be moved, it is sometimes possible to track the relative linear and angular motion. For example, in state of the art games for the Nintendo Wii, players are instructed on how to hold and move their controller. By assuming the players are approximately following the instructions, it is possible to interpret the signal from the tri-axial accelerometer to roughly track the relative linear and angular motion of the controller. But there is a wide class of games and other applications where it is undesirable to constrain how the user may move or hold the movable controller. State of the art movable controllers are therefore unnecessarily limiting.\nTo review the background of sensing the positions and tracking the paths of objects moving in a three dimensional space, the tracking is done by inertial navigation systems (INS) which use a combination of accelerometers and gyroscopes to create or compute an inertial frame within which accelerations represent strictly linear acceleration in the world frame. If one knows the world frame linear acceleration of an object over time, one can calculate the current position of that object over time with respect to its starting location. If one knows the angular velocities of an object over time, one can provide its orientation at any point in time. Conventionally, in the tracking of objects, linear accelerations combined with angular velocities are necessary and sufficient for providing location and orientation of an object with respect to a starting location. There are six unknowns that must be solved for at every point in time. Most INSs (inertial navigation systems) use gyroscopes to fix or solve for the three angular velocities. Once the orientation over time is known, accelerometers can be used to track the three linear accelerations as described above. Reference is made to the publication, “Radar, Sonar, Navigation & Avionics Strapdown Inertial Navigation Technology, 2nd Edition”, by D. Titterton and J. Weston, published in 2005 as part of the IEE Radar, Sonar, Navigation and Avionics Series, for an introduction to and further information on the field of inertial navigation systems.\nReference is also made to the publication, “Design and Error Analysis of Accelerometer-Based Inertial Navigation Systems”, Chin-Woo Tan et al., Published in June, 2002 by the University of California at Berkeley for the State of California PATH Transit and Highway System which is hereby incorporated by reference. This is a study of the feasibility of inertial navigation systems that use only accelerometers to compute the linear and angular motions of a rigid body. This publication relates to systems for measuring linear and angular velocities of motor vehicles and the like. Its determinations track motion on scale of tens of meters accuracy on a time scale of tens of minutes."}
{"text": "A user interface (abbreviated to U.I.) is the combination of at least a means of receiving commands and at least a means of display. A U.I. generally uses descriptive data for preparing menus capable of being displayed on a television screen. Amongst other things, these data can be used to prepare lists of items that the user may select. The MPEG-7 standard specifies descriptive data associated with audiovisual contents. The MPEG-7 standard can be accessed via the ISO under the reference ISO/IEC JTC1/SC29/WG11-N5525 published in March 2003. If these data originate from a transmission network, they can be stored in a local database and continuously updated according to the contents transmitted by the network. There are other descriptive data, also read from media such as CDs or DVDs. Generally speaking, the descriptive data are recorded in a receiver memory database and grouped together by lists according to certain criteria such as: the titles of the programs transmitted, the transmission channels, the program topics, the available languages, etc. Other lists may also be prepared by the receiver itself, lists of functions for example.\nSome lists are accessible with the aid of an Electronic Program Guide (abbreviated to EPG) displayed on a screen. The EPG retrieves list items from a transmission network for example and is responsible for displaying an item identifier in a box. The boxes are arranged vertically or horizontally, or in a grid. The user navigates the displayed lists with the aid of navigation keys arranged on a remote control. If the boxes are arranged horizontally, then the user navigates the list using the “right arrow” or “left arrow” keys. When a box is highlighted graphically, the item associated with this box is selected. Other keys can be used to execute functions applied to this item, for example if the item is an audiovisual program, the functions that can be executed may be viewing it, recording it or associating it with a text for example, etc. In some EPGs, the boxes associated with selectable items are spread over the screen, either within a structured list, or in isolation. The user may have difficulties in quickly selecting the item that they want since starting from a given item, the choice of navigation key to use is not clearly apparent to them.\nThe document WO01/43424—UNITED VIDEO PROPERTIES discloses the menu display forming lists arranged vertically. At the beginning and end of the list, icons accessible via navigation controls indicate that there are other items that would be displayed by continuing to navigate in this direction. For example, FIG. 28 and the corresponding text show the display of a favorite list, the display of a favorite list is triggered by selecting an icon in the icon field of FIG. 16. In FIG. 28, the icon that appears above an item list enables a sort function to be applied to the list below this icon. This document teaches access to a single set of functions from an additional icon.\nThe list displayed may comprise a large number of items. A well-known solution consists in using two page jumps forward and back keys. But this means having to increase the number of keys, which increases the cost of the remote control.\nIn addition to the functionalities that it provides, the remote control forms an important aesthetic element in an audiovisual system. According to fashion, they comprise more or less keys, and these keys take various forms. It may be useful to produce new forms of remote control with uncluttered lines and whose number or form of the keys arouses the user's interest.\nThe present invention discloses a new way of selecting an item within a list or a function, with limited and inexpensive means of receiving commands."}
{"text": "1. Field of the Invention\nThe present invention relates to an image pickup unit, mounted in a facsimile or a cash dispenser, for taking an image of the surface of an object (more specifically, printed characters on an original or embossed characters on a cash card, for example).\n2. Prior Art\nFIG. 2 is a sectional view showing an example of a conventional specular reflection type image pickup unit having a image-forming means, such as a rod lens, arranged at such a position as to receive the maximum quantity of reflected light.\nIn this image pickup unit, a light source 1 is provided at a position a specified distance away from the target surface of the object, such as an original P. The light source 1 is formed of a plurality of light emitters arranged in a row, and emits light ph1 in a linear form at a specified angle xcex8(0xc2x0 less than xcex8 less than 90xc2x0) to a direction normal to the target surface of the original P. After passing through a glass plate 2, the emitted light ph1 is incident on the focal points f of a rod lens array 3 on the target surface of the original P at an incidence angle xcex8. The incident light is reflected by the target surface chiefly as a maximum quantity of reflected light rf1 at a reflection angle xcex8. For this purpose, the rod lens array 3 is so arranged as to hold the angle xcex8. Therefore, when the reflected light rf1 is converged by the rod lens array 3, an image is formed at an image-forming position indicated by a conjugate distance T. Note that the conjugate distance T is the characteristic value of the rod lens array 3 that shows the distance between the target surface of the original P and the image-forming position. The light-receiving faces of image pickup elements 4 are located at this image-forming position. The reflected light rf1 which formed the image plane is converted into electric signals by an array of image pickup elements 4, and the electric signals pass through a board 4a and are sent to an image processing system, for example, to undergo a specified process. The light source 1 is mounted on the board 1a, and the board 1a is located at a proper position and supported in a housing 5. The image pickup elements 4 are fixed through the intermediary of a board 4a to the inside walls of the housing 5 such that the light-receiving faces are normal to the propagating direction of the reflected light rf1.\nAs has been described, in the specular reflection type image pickup unit, the light source 1 and the rod lens array 3 are arranged at equal angles xcex8 to the normal to the target surface of the original P. Therefore, since the image pickup elements can receive the maximum quantity of the reflected light, the embossed characters of a cash card, for example, can be displayed with proper contrast on the display screen.\nHowever, the conventional image pickup unit in FIG. 2 has a problem. Specifically, the problem is the large width of the housing attributable to the inclined mounting in the housing 5 of the board 4a on which the image pickup elements 4 are mounted.\nMeanwhile, FIG. 3 is a sectional view showing an example of the irregular reflection type image pickup unit whose size was reduced by the sacrifice of the received light quantity of the image-forming means. The same parts as those in FIG. 2 are designated by the common reference numerals.\nIn this image pickup unit in FIG. 3, the rod lens array 3 is arranged normal to the target surface of the original P, and receives the reflected light rf1. Therefore, the width dimension of the housing 5A can be made smaller than that of the housing 5 in FIG. 2. However, due to the smaller received quantity of the reflected light rf1, this image pickup unit is unable to produce a distinct image.\nIt is an object of the present invention to provide an image pickup unit which is made in a compact size but capable of producing vivid images.\nIn present invention, the image pickup unit comprises a light source, an image-forming means, image pickup elements, a housing, and a board.\nThe light source for emits a light to a target surface of an object at a fixed angle to a direction normal to the target surface. The image-forming means is mounted at such an angle as to receive a maximum quantity of reflected light from the target surface, and forms an image at a specified image-forming position. The image pickup elements receive the reflected light converged at the image-forming position and convert the reflected light into an electric signal. The light source, the image-forming means, and the image pickup elements are arranged in the housing. The image pickup elements are mounted on the board, and the board is arranged in the housing in parallel with the target surface.\nBy adopting the configuration described above, the space between two walls of the housing which are parallel with the target surface of the object can be reduced, so that the housing itself can be reduced in size.\nBecause the image-forming means is arranged at such an angle as to receive a maximum quantity of the reflected light, distinct images can be obtained by the image pickup elements.\nAccording to a second embodiment of the present invention, an image pickup unit comprises a light source, an image-forming means, image pickup elements, a housing and a board.\nThe light source emits a light to a target surface of an object at a fixed angle to a direction normal to the target surface. The image-forming means is mounted at such an angle as to receive a maximum quantity of reflected light from the target surface, and forms an image at a specified image-forming position. The image pickup elements receive the reflected light converged at the image-forming position and convert the reflected light into an electric signal. The light source, the image-forming means, and the image pickup elements are arranged in the housing. The image pickup elements are mounted on the board, and the board is arranged in the housing in a direction at right angles to the target surface.\nBy adopting the configuration shown in the second embodiment, the space between two walls of the housing, each of which lies in a direction at right angles to the target surface of the object, can be reduced, with the result that the housing itself can be reduced in size.\nLike in the first embodiment, because the image-forming means is arranged at such an angle as to receive a maximum quantity of the reflected light, vivid images can be obtained by the image pickup elements.\nIn the first and the second embodiments, that portion of the board where the image pickup elements are mounted may be bent to orient the light-receiving faces of the image pickup elements at right angles to the propagating direction of the reflected light. By this arrangement, the image pickup elements can be positioned correctly at the image-forming means merely by moving the board in one direction along the wall.\nIn the first and the second inventions, it is possible to arrange another board for mounting the light source in a direction at right angles to the target surface and install a reflector to reflect the emitted light of the light source at a fixed angle. By this arrangement, the other board, which was inclined, is now arranged in parallel with the wall opposite the wall on which the other board is mounted, so that the housing can be further reduced in size."}
{"text": "1. Field of the Invention\nThe present invention is directed, in general, to power supplies and, more specifically, to a power supply embodying a volume-efficient construction allowing the power supply to fit within the chassis of a portable computer system.\n2. Description of Related Art\nPortable, battery-powered computers have become increasingly popular over the last several years due to their light weight and small size that permit them to be easily hand-carried in an ordinary briefcase and used by business travelers in cramped spaces, such as on airline seat back trays, lacking electrical plug-in facilities. In fact, a particularly small type of portable computer, the notebook computer, is very popular, generally having dimensions of 8.5\".times.11\" and a weight of less than 8 pounds. More recent developments in computer miniaturization have resulted in so-called \"subnotebook\" computers having still smaller dimensions.\nThe modern portable computer typically incorporates both hard and floppy disk drives, a display screen built into a pivoting display screen portion, and a keyboard built into a main chassis portion. It is thus a fully self-contained computer able to be conventionally used, for at least short periods of time, in situations and locations in which the use of a much larger desktop computer is simply not feasible.\nAs is well-known, however, even state-of-the-art portable computers have certain limitations and disadvantages, compared to their much larger desktop computer counterparts. One principal disadvantage is that portable computers are battery-powered and thus are subject to the limited life of the battery. Fortunately, however, almost all battery-powered portable computers are sold with an AC adapter. When plugged into a convenient source of AC power, the AC adapter converts the power into a form usable by the portable computer. The AC adapter comprises a line cord that plugs into the AC power source and an adapter cord that plugs into a designated power supply port usually on a backside of the portable computer main chassis. Power conversion is performed in a power supply that exists in the form of a power supply chassis separate from that of the portable computer and coupled between the line cord and adapter cord.\nThere are historical reasons for the power supply being housed in a separate power supply chassis. Because of notoriously tight volumetric constraints in portable computers main chassis, the components and subsystems therein have for years been the subject of constant, extensive design work directed toward making the components and subsystems as small as possible. Because the main chassis has always been fully occupied with necessary components and subsystems, there was no motivation to attempt to fit in less vital subsystems, such as the power supply. Thus, because prior art power supplies have been entirely external to the portable computer, they were never subjected to the intense miniaturization effort otherwise universally applied.\nAccordingly, prior art separate supplies have been comparatively large and over-engineered. Heat-producing components in such supplies are most often cooled by mounting heat sinks to each component. The fins or projections of the heat sinks are bulky and heavy. The power supplies are surrounded by an electromagnetic interference (\"EMI\") shield to attenuate and guard against stray EMI. Thus, the prior art power supply chassis housing these power supplies and their associated shields have been correspondingly large, heavy and cumbersome, earning them the unflattering colloquial term \"brick.\" The bulky and clumsy \"brick\" stood in stark contrast to the otherwise elegant and compact design of prior art portable computers.\nAgain, due to the tight volumetric requirements of components in the portable computer, it has not as yet been feasible to relocate the brick to within the portable computer main chassis. The power supply must undergo the same miniaturization effort through which the rest of the portable computer has already undergone. However, miniaturization of a power supply is not trivial. It not only entails an electrical redesign effort to minimize the size of the electronic components thereof, but also requires a fundamental rethinking of related mechanical, electromagnetic and thermal issues.\nFirst, a volumetrically efficient means of mechanically attaching the power supply chassis within the portable computer chassis must be developed. Second, the power supply must be electromagnetically shielded from other components and subsystems in the portable computer, again, in a space-saving manner. Finally, and perhaps most critically, new techniques for removing heat from selected power supply components must be developed, as it is no longer acceptable to consume space with the individual finned heat sinks so often resorted to in the prior art.\nAccordingly, what is needed in the art is a volume-efficient construction for a portable computer power supply that allows the power supply to fit within the confines of a portable computer system."}
{"text": "1. Field of the invention\nThis invention relates to a method of producing palatinose for sucrose using immobilized .alpha.-glucosyl transferase.\n2. Description of the Prior Art\nPalatinose is a reducing disaccharide and it has been confirmed that palatinose is a natural component in honey and cane juice. The sweetness of it is about half of that of sucrose.\nThe molecular structure of palatinose is 6-O-.alpha.-D-glucopyranosyl-D-fructose in which glucose is bonded to fructose at 1,6-position. Palatinose is also called as isomaltulose.\nIn regard to a method of producing palatinose, a method is known wherein enzyme in a form of living or dead cells of Protaminobacter rubrum is added to an aqueous solution of sucrose to convert sucrose into palatinose and then the enzyme is separated from the solution as disclosed in West German Pat. No. 1,049,800. In this case, the separation of the enzyme is performed by centrifugal separation.\nAccording to the technical level, .alpha.-glucosyl transferase forming palatinose from sucrose is produced by cultivating a certain kind of bacteria, i.e., Protaminobacter rubrum, Serratis plymuthica, etc., in the presence of sucrose. Since the enzyme associates with the cell of the bacteria, a mass of their cells itself can be used as an enzyme preparation.\nTherefore, in the case of producing palatinose from sucrose, a mass of cells containing .alpha.-glucosyl transferase is added to an aqueous solution of sucrose in a concentration of 20-30% and the reaction is performed by stirring the mixture while controlling the mixture at a temperature of 20.degree.-30.degree. C. and a pH of about 7. The reaction period of time depends upon the amount of enzyme, the reaction temperature, etc., but, since the remaining amount of sucrose in the reaction mixture becomes very slight usually in about 20 hours after the addition of the cells, this point is defined as the end of the reaction, and then the cells are removed from the reaction mixture to provide an aqueous transparent solution of palatinose, which is purified by ion-exchange resin, etc., and then concentrated under reduced pressure to crystallize palatinose.\nIn the above procedure, the cells are removed from the reaction mixture usually by a centrifugal separation method, but, since the difference in specific gravity between the cells and the reaction mixture is small and also the size of the cells is very small, the separation efficiency is very poor and hence a very large installation is required for practicing the method, which gives a large problem in the case of practicing the method industrially."}
{"text": "For many years, there has been considerable research in optical digital computing systems with the objective of developing and implementing parallel processing systems. The advantages of using optical systems allowing high throughput parallel processing with very high space-bandwidth products are well known to the persons skilled in the art. Nevertheless, most optical systems heretofore used have been analog in nature and have therefore sacrificed some of the accuracy and flexibility that digital computing can provide. Furthermore, analog optical systems have severe limitations in the sizes of the operational gates. Continuing efforts in which the object of the present invention partakes have been made to combine optical systems and digital algorithms in order to obtain the benefits of both.\nIn the article entitled \"Coherent optical implementation of generalized two-dimensional transforms\", by James R. Leger and Sing. H. Lee, Optical Engineering, Vol. 18, no. 5, 1979, there is disclosed a coded phase coherent optical processor capable of performing two dimensional linear transformations. The set of two-dimensional Walsh functions is chosen as a transform basis. The coherent optical correlator performs a correlation function between a two-dimensional object image and the Walsh-Hadamard transform of that image. The optical output plane consists of delta functions which represent the maximum of eight discrete levels of the Walsh-Hadamard transformation. However, without the accurate computation of spectral expansion functions and accurate data encoding of computer generated holograms, the coded phase optical correlator can only play a limited role for optical image processing and optical pattern recognition. This article is hereby incorporated by reference.\nIn U.S. Pat. No. 4,318,581 to Guest, published in March 1982 and in the article entitled \"Two proposed holographic numerical optical processors\", by C. C. Guest and T. K. Gaylord, SPIE, Vol. 185, 1979, there are disclosed two holographic numerical optical processors based on EXCLUSIVE OR processing and NAND-OR-OR processing. Input data words of the truth table perform either EXCLUSIVE OR or NAND-OR-OR processing on partially recorded truth table input data word array with truth values \"1\". After logical operations have occurred, the output row photoconductor arrays are filled with the necessary patterns such that the truth table output data words can be reconstructed through the photoconductor arrays. The optical processors described in the two references hereabove mentioned, which are incorporated by reference, provided an advance in the prior art of the time insofar as they had a simple structure. Furthermore, the mapping relation existed for both the location addressable memory and the content addressable memory. Additionally, the truth table look up processing could be performed for either numerical arithmetic operations or boolean logic operations. However, in the Guest patent, the digital inputs and digital outputs are one-dimensional. There is no teaching in the Guest patent of any possibility to extend the computational methods to a two-dimensional input. Futhermore, the hologram recorded in the storage medium comprises solely the pattern of the corresponding truth table input pattern corresponding to the value 1 as an output value. Such hologram recording saves considerable space but can only allow logic operations on the digital inputs. Guest illustrates it with the use of the XOR and the NAND-OR-OR logical operations. The XOR operation, for instance, is performed by appropriately phase shifting between the reference beam and the object beam. When the number of digital inputs becomes large, severe problems arise in the recording of large truth-tables in the hologram, thereby limiting the use of the Guest processor to modest input sizes.\nAnother recent approach to optical truth table look-up processing has been disclosed in an article by Shing-Hon Lin, Thomas F. Krile and John F. Walkup, entitled \"Optical triple-product processing in logic design\", Applied Optics, Vol 25, no. 18, 1986, which is hereby incorporated by reference. In this article, triple-product processing is used to perform a generalized bilinear transform on two different input functions to produce the outputs of the truth table. Yet, the triple-product processor utilizes two input planes, thereby preventing the two-dimensional space bandwidth product from being fully used. An additional drawback resides in the lack of accuracy in the computation of the spectral mapping coefficients by the generalized bilinear transform for truth table mapping.\nAll the disadvantages hereabove described in connection with optical systems of the prior art are overcome by the optical system of the present invention. The system architecture according to the present invention is intrinsically a massively parallel logic device which can run at the highest possible speed permitted in operations of basic micro-level computing functions. The optical system of the present invention is equivalent to an optical two-dimensional, large-scale integrated circuit with at most a one- or two-gate delay. Not only can it run at extremely high speeds, but it obviates the problems caused by radiation and the spurious effects produced by electromagnetic field interferences."}
{"text": "1. Field of the Invention\nThis invention relates to the biological modification of the lignocellulosic components of gramineous agricultural plant materials for the purpose of making the cellulose more available for use in a variety of applications. For instance, the digestibility of lignocellulosic grasses by cattle is limited to only that portion of the cellulose which is not physically bound by the lignin component. Also, with the current trend toward the production of alcohol as a renewable energy resource, the feasibility of employing lignocellulosic residues as sources for fermentable sugars hinges on the amount of free cellulose which can be enzymatically hydrolyzed to a glucose. Typically, up to about 50% of the cellulose in natural, grassy plant tissue is rendered unavailable for such uses by the lignin associated therewith.\n2. Description of the Prior Art\nStudies of the degradation of lignocellulosic materials have revealed that the lignin component can be depleted disproportionately faster than the polysaccharides when acted on by some white-rot fungi [Kirk, T. K., Phytopathology 63: 1504-1507 (1973)]. An approach to optimizing the ability of such organisms to preferentially decompose the lignin in wood materials is taught by Eriksson et al., U.S. Pat. No. 3,962,033 and P. Ander et al., Material and Organismen, Beiheft 3: 129-140, which apply a procedure for mutating spores of mycelia of wild strain white-rot fungi, and selecting the mutants having the desired degradative properties. While Eriksson teaches that these organisms can also be used for reducing the lignin content of certain grassy materials (e.g., straw and bagasse), their extent of selectivity in degrading the components of these materials is unreported. The chemical structure of grassy lignins is characteristically different from that of the woody lignins, and an organism adapted for preferentially degrading lignocellulose in one, typically will not manifest the same degree of preferentiality in the other.\nZadrazil [European J. Appl. Microbiol. 4: 273-281 (1977)] observed that a variety of basidiomycetes, including Pleurotus cornucopiae, Pleurotus sp. (Florida isolate), and Stropharia rugosannulata would generally increase the digestibility of wheat straw after about 20 days of cultivation. Kaneshiro [Dev. Ind. Microbiol. 18: 591-598 (1977)] reports preliminary results that the basidiomycetes Pleurotus ostreatus has potential for preferentially degrading the lignin content of feedlot waste fiber and wheat straw. In one fermentation experiment, the lignin content in a mixture of the fiber and straw was decreased from 12% to 7% with only slight depletion of the cellulose."}
{"text": "1. Field of the Invention\nThis invention relates to a recording sheet for ink jet printer (hereinafter refered to as ink jet recording sheet).\n2. Prior Art\nIn recent years, the demand for color printer is increased. Particularly, an ink jet printer, one of the non-impact recording systems, is highly increased since a comparatively rapid color recording is possible for its simplified system, since the noise is less than impact printers such as dot-printer, etc., and since the use in offices is easy, and the like.\nIn ink jet recording sheets, the use of synthetic amorphous silica has been proposed in various forms for the purpose of improving the color density, ink absorbency, ink drying ability, feathering of ink and flowing out of ink. For example, there are disclosed an ink jet recording sheet coated with a non-colloidal silica powder in Japanese Laid-Open Patent Application No. 55-51583, an ink jet recording sheet having a base paper coated with a mixture of a fine silicic acid and a water-soluble resin in Japanese Laid-Open Patent Application No. 56-148583, and an ink jet recording sheet obtained by at least double applications of the same coating color consisting of both synthetic silica and aqueous binder in total coating-amount of at least 10 g/m.sup.2 in Japanese Laid-Open Patent Application No. 57-107879.\nAs synthetic amorphous silica, there are silica manufactured by a wet precipitation process, by a wet gel process or by a humed process, i.e. a process of hydrolysing silicium tetrachloride in a oxyhydrogen flame. In conventional ink jet recording papers, silica manufactured by a wet precipitation process or by a wet gel process is mainly used, and it is disclosed in the above Laid-Open Patent Applications.\nThe use of silica manufactured by a humed process is disclosed in Laid-Open Patent Application 60-204390, which describes that an ink receptive layer comprises as the feature a synthetic ultra-fine silica manufactured by a vapor phase process. However, this Patent Application practically discloses a method for coating a ultra-fine silica on an ink receptive layer which is obtained by coating a synthetic amorphous silica in a coating amount of at least 10 g/m.sup.2 on a substrate, and it discloses a method for coating a mixture of an ultra-fine silica and a silica manufactured by a precipitation process in a certain mixed-ratio on a substrate in a coating amount of 15 g/m.sup.2. In any cases, this Laid-Open Patent Application discloses a coating amount of at least 10 g/m.sup.2. Accordingly, the technology in Japanese Laid-Open Patent Application No. 60-204390 is substantially equal to the technology disclosed in Japanese Laid-Open Patent Application Nos. 55-51583, 56-148583, 57-107879 and the like.\nThat is, the purpose of the above Japanese Laid-Open Patent Publications consists in the preparation of an ink jet recording sheet having both a good recording density and an excellent ink absorbency by applying a coating material of 10-20 g/m.sup.2 on a substrate.\nAs methods for preparing an ink jet recording sheet having an improved water-resistance, Japanese Laid-Open Patent Application No. 56-84992 discloses the coating of a polycation-polyelectiolyte on the surface of an ink receptive layer, and Japanese Laid-Open Patent Application No. 59-20696 discloses the use of an ink receptive layer containing polydimethyldiallyl ammonium chloride.\nAs the method for preparing a recording sheet having both a good light resistance and an excellent image-light resistance, there is wellknown various ink receptive layers containing each of metal oxide, UV-absorber, oxidation inhibitor, etc.\nJapanese Laid-Open Patent Application No. 58-177390 discloses that an ink jet recording sheet suitable for an ink jet printer of electric field control type is obtained by coating an electric conductive agent of quaternary ammonium type on the surface of an ink receptive layer containing both a synthetic silica and an aqueous binder for the purpose of providing a normal printability, wherein the sheet is independent of environmental humiditic changes by means of preventing the electrostatic chargeability.\nAs described above, the subject of the hitherto ink jet recording sheets was to improve various properties, that is, the color density, the ink-absorbency, drying ability, feathering, flowing out, the water resistance and light resistance of a image and a recording sheet (printability of a recording sheet), and the like. Recently, a discoloration and a fading caused by indoor and outdoor gases have become a great problem.\nAs gases causing a discoloration or a fading, there are NO.sub.x -gas, sulfur dioxide gas, etc., in exhaust gases or factory smokes; ozone generated from many copying machines; smokes of cigarets; and the like.\nOzone is generated more and more, since many small or medium copying machines have been used increasingly in offices owing to the miniaturization and the simplification of many copying machines through the Office Automations. However, the discoloration or the fading caused by indoor and outdoor gases (oxidizing gases) forms a particularly great problem. It is wellknown that some of oxidizing gases cause an oxidizing excitation of acidic dye-molecules in an aqueous ink for an ink jet recording.\nFor the preparation of a sheet having both an appropriate recording density and a sufficient ink absorbency in an full color ink jet recording with high quality and density, it is required that an ink-receptive layer is provided in accordance with the maximum amount of ejected inks in an unit printing area. In this case, the conventional technique of coating a conventional synthetic amorphous silica such as a silica manufactured by a wet process or by a humed process, usually requires a coating amount of at least 10 g/m.sup.2 on a substrate.\nIn coating a filler such as silica etc. on a fiber-like substrate in an amount of circa 10 g/m.sup.2, the fibers are covered with powdered fillers completely, so that the paper-like figure and feeling are lost in a paper-base substrate. And owing to the less bulk density, the fine silica requires more binder than other fillers.\nSince many vacant spaces in the coating layer are filled with the binder, the coating amount must be further increased for providing a good ink-absorbency, so that the paper-like figure and feeling are increasingly lost. In a small amount of binder, the recording sheet has a good ink absorbency, but it exhibits as the defects a weak surface strength and the troubles of paper-powder generation, of clogging the oriffices by the powders and the like, so that it is not used as a usual recording sheet.\nThe synthetic silica, which is used as a catalyst in oxidizing-reactions, has an excellent function as a catalyst. Accordingly, in the conventional recording sheets coated with a large amount of synthetic silica for improving both a recording density and an ink absorbency, there are following defects: With an increased amount of synthetic silica, the fading of the image is strengthen by oxidizing gases generated in the catalytic action, and the paper-like figure and feeling are further decreased.\nMoreover, metal oxide, UV-absorber, oxidation-inhibitator, polycation-polyelectrolyte. etc., which are used for improving both a light resistance of a recording sheet and that of an image, have no ability of preventing a fading caused by oxidizing gases. And some of the above agents increase the fading.\nAs described above, there is unknown a method for preventing the fading of the ink jet recording-image by oxidizing gases, etc., under maintaining both a good recording density and a high ink absorbency with excellent paper-like figure and feeling.\nFurther, it has been found that such a method is not achieved by extending the conventional techniques."}
{"text": "The present invention relates to the microwave generation, amplification, and processing arts. It particularly relates to traveling wave tubes for microwave amplifiers and microwave oscillators, and will be described with particular reference thereto. However, the invention will also find application in other devices that operate at microwave frequencies, and in other devices that employ slow wave circuits.\nTraveling wave tubes typically include a slow wave circuit defined by a generally hollow vacuum-tight barrel with optional additional microwave circuitry disposed inside the barrel. An electron source and suitable steering magnets or electric fields are arranged around the slow wave circuit to pass an electron beam through the generally hollow beam tunnel. The electrons interact with the slow wave circuit, and energy of the electron beam is transferred into microwaves that are guided by the slow wave circuit. Such traveling wave tubes provide microwave generation and microwave amplification.\nHeretofore, commercially produced traveling wave tubes have been limited to about 65 GHz. However, future applications call for traveling wave tubes that operate at frequencies of 100 GHz or higher. For space-based applications these high frequency devices will probably be driven at operating voltages of 20 kV or less in accord with presently available space-based electrical power sources.\nConstruction of high frequency traveling wave tubes is difficult using existing traveling wave tube manufacturing techniques. Designs for high frequencies call for microwave circuitry with very small features (for example, repetition periods of less than 0.2 mm), and greatly reduced quantities of dielectric insulation material in the tube to reduce dielectric loading. Moreover, adequate heat sinking becomes an increasingly significant issue as the operating frequency increases.\nIn one known method for manufacturing traveling wave tubes, a fragile three-dimensional microwave circuit shell, such as a metallic helix, is compressively secured in a generally hollow cylindrical barrel. Dielectric rods arranged parallel to the helical axis of the microwave circuit act as standoff insulators that align and secure the compressed microwave circuit shell inside the barrel.\nTo ensure good thermal contact between the components, the compressive forces are substantial. The fragile microwave circuit shell and dielectric rods are compressively secured in the barrel by briefly heating the barrel during insertion to induce temporary thermal expansion of the barrel. The microwave circuit shell/dielectric rods combination has close tolerances with respect to the barrel, and so when the barrel contracts upon cooling the interior components are compressively secured in the barrel. However, the heating and compression can damage the slow wave circuit, and mass production by this method is difficult. Moreover, this technique is not well suited for the small structures used in devices that are preferred for 100 GHz or higher operation.\nTo achieve features on the delicate scale called for in high frequency operation, lithographic techniques are regularly used in the electronics industry. However, these techniques are generally applied to planar wafer substrates of silicon or other semiconductor materials. Lithographic techniques are not readily adapted to produce the types of finely detailed three-dimensional structures called for in traveling wave tubes designed for high frequency operation.\nTo reduce dielectric loading, the dielectric rods can be replaced by thin dielectric standoff chips of natural diamond. In one constructed device described in A Millimeter-Wave TunneLadder TWT (D. Wilson, NASA Contract Report 182183, 1988), diamond chips with heights of 250 microns each were used in a traveling wave tube that operated at 28 GHz. However, this device has not been replicated to date due to the cost of natural diamond and the assembly difficulties, especially relating to positioning of the diamond chips. Moreover, scaling such a device up to 100 GHz or higher frequency would call for a large number of diamond chips (e.g., about 80-150 diamond chips for a 2-3 cm long traveling wave tube designed for 100 GHz operation) each having a height of about 75 microns or thinner. These reduced dimensions and increased numbers of diamond chips would further exacerbate an already difficult manufacturing process.\nThe present invention provides an improved apparatus and method."}
{"text": "Glaucoma is an ocular disorder associated with elevated intraocular pressures which are too high for normal function and may result in irreversible loss of visual function. If untreated, glaucoma may eventually lead to blindness. Ocular hypertension, i.e., the condition of elevated intraocular pressure without optic nerve head damage or characteristic glaucomatous visual field defects, is now believed by many ophthalmologists to represent the earliest phase of glaucoma.\nMany of the drugs formerly used to treat glaucoma proved not entirely satisfactory. Indeed, few advances were made in the treatment of glaucoma since pilocarpine and physostigmine were introduced. Only recently have clinicians noted that many .beta.-adrenergic blocking agents are effective in reducing intraocular pressure. While many of these agents are effective in reducing intraocular pressure, they also have other characteristics, e.g. membrane stabilizing activity, that are not acceptable for chronic ocular use. (S)-1-tert-Butylamino-[(4-morpholino-1,2,5-thiadiazol- 3-yl)oxy]-2-propanol, a .beta.-adrenergic blocking agent, was found to reduce intraocular pressure and to be devoid of many unwanted side effects associated with pilocarpine and, in addition, to possess advantages over many other .beta.-adrenergic blocking agents, e.g. to be devoid of local anesthetic properties, to have a long duration of activity, and to display minimal tolerance.\nAlthough pilocarpine, physostigmine and the .beta.-blocking agents mentioned above reduce intraocular pressure, none of these drugs manifests its action by inhibiting the enzyme carbonic anhydrase and, thereby, impeding the contribution to aqueous humor formation made by the carbonic anhydrase pathway.\nAgents referred to as carbonic anhydrase inhibitors block or impede this inflow pathway by inhibiting the enzyme, carbonic anhydrase. While such carbonic anhydrase inhibitors are now used to treat intraocular pressure by oral, intravenous or other systemic routes, they thereby have the distinct disadvantage of inhibiting carbonic anhydrase throughout the entire body. Such a gross disruption of a basic enzyme system is justified only during an acute attack of alarmingly elevated intraocular pressure, or when no other agent is effective. Despite the desirability of directing the carbonic anhydrase inhibitor only to the desired ophthalmic target tissue, no topically effective carbonic anhydrase inhibitors are available for clinical use.\nHowever, topically effective carbonic anhydrase inhibitors are reported in U.S. Pat. Nos. 4,386,098; 4,416,890; and 4,426,388. The compounds reported therein are 5 (and 6)-hydroxy-2-benzothiazolesulfonamides and acyl esters thereof."}
{"text": "This invention relates to digital data processing systems and, more particularly, to multiprocessing systems with distributed hierarchical memory architectures.\nThe art provides a number of configurations for coupling the processing units of multiprocessing systems. Among the earlier designs, processing units that shared data stored in system memory banks were coupled to those banks via high-bandwidth shared buses or switching networks. During periods of heavy usage, bottlenecks were likely to develop as multiple processing units simultaneously contended for access to the shared data.\nIn order to minimize the risk of formation of transmission bottlenecks, distributed memory systems were developed coupling individual processing units with local memory elements to form semi-autonomous processing cells. To achieve the benefits of multiprocessing, some of the more recently designed systems established cell communications through utilization of hierarchical architectures.\nThe distributed memory systems, however, permit multiple copies of single data items to reside within multiple processing cells; hence, it is difficult insure that all processing cells maintain identical copies of like data elements. Conventional efforts to resolve this problem, i.e., to preserve data coherency, rely upon software oriented techniques utilizing complex signalling mechanisms.\nTo avoid processing and signalling overhead associated with these software oriented solutions, Frank et al, U.S. Pat. No. 4,622,631, discloses a multiprocessing system in which a plurality of processors, each having an associated private memory, or cache, share data contained in a main memory element. Data within that common memory is partitioned into blocks, each of which can be owned by any one of the main memory and the plural processors. The current owner of a data block is said to have the correct data for that block.\nA hierarchical approach is disclosed by Wilson Jr. et al, United Kingdom Patent Application No. 2,178,205, wherein a multiprocessing system is said to include distributed cache memory elements coupled with one another over a first bus. A second, higher level cache memory, attached to the first bus and to either a still higher level cache or to the main system memory, retains copies of every memory location in the caches below it. The still higher level caches, if any, and system main memory, in turn, retain copies of each memory location of cache below them. The Wilson Jr. et al processors are understood to transmit modified copies of data from their own dedicated caches to associated higher level caches and to the system main memory, while concurrently signalling other caches to invalidate their own copies of that newly-modified data.\nNotwithstanding the solutions proposed by Frank et al and Wilson Jr. et al proposal, data coherency and bus contention remain significant problems facing both designers and users of multiprocessing systems. With respect to Wilson Jr. et al, for example, these problems may be attributed, at least in part, to the requirement that data in main memory must always be updated to reflect permanent modifications introduced to the data elements by each of the processors in the system. Moreover, neither of the proposed designs is capable of supporting more than a limited number of processing units. This restriction in xe2x80x9cscalabilityxe2x80x9d arises from a requirement of both the Wilson Jr. et al and Frank et al systems that the size of main memory must increase to accommodate each additional processor.\nIt is therefore an object of this invention to provide an improved multiprocessing system with improved data coherency, as well as reduced latency and bus contention. A further object is to provide a multiprocessing system with unlimited scalability.\nOther objects of the invention are to provide a physically distributed memory multiprocessing system which requires little or no software overhead to maintain data coherency, as well as to provide a multiprocessing system with increased bus bandwidth and improved synchronization.\nThe invention attains the aforementioned objects by providing, in one broad aspect, a digital data processing system comprising a plurality of processing cells arranged in a hierarchy of rings. The system selectively allocates storage and moves exclusive data copies from cell to cell in response to access requests generated by the cells. Routing elements are employed to selectively broadcast data access requests, updates and transfers on the rings.\nA system of the type provided by the invention does not require a main memory element, i.e., a memory element coupled to and shared by the systems many processors. Rather, data maintained by the system is distributed, both on exclusive and shared bases, among the memory elements associated with those processors. Modifications to datum stored exclusively in any one processing cell do not have to be communicated along the bus structure to other storage areas. As a result of this design, only that data which the processors dynamically share, e.g., sharing required by the executing program themselves, must be transmitted along the bus structure. These aspects, along with the systems hierarchical structure, localize signal traffic greatly, thereby reducing bus contention and bottlenecks.\nWith further attention to system structure and element interconnection, the processing cells include central processing units coupled with memory elements, each including a physical data and control signal store, a directory, and a control element. Groups of cells are interconnected along unidirectional intercellular bus rings, forming units referred to as segments. These segments together form a larger unit referred to as xe2x80x9cinformation transfer domain(0).xe2x80x9d While cells residing within each segment may communicate directly with one another via the associated intercellular bus, the associated central processing units are not themselves interconnected. Rather, intersegment communications are carried out via the exchange of data and control signals stored in the memory elements. A memory management element facilitates this transfer of information.\nCommunications between cells of different domain(0) segments are carried out on higher level information transfer domains. These higher level domains are made up of one or more segments, each comprising a plurality of domain routing elements coupled via a unidirectional bus ring. It will be appreciated that the segments of higher level domains differ from those of domain(0) insofar as the former comprise a ring of routing elements, while the latter comprise a ring of processing cells. Each routing element is connected with an associated one of the segments of the next lower information transfer domain. These connected lower segments are referred to as xe2x80x9cdescendants.xe2x80x9d Every information transfer domain includes fewer segments than the next lower domain. Apart from the single segment of the system\"\"s highest level domain, signals are transferred between segments of each information transfer domain via segments of the next higher domain.\nAn exemplary system having six domain(0) segments includes two domain(1) segments, the first which transfers data between a first three of the domain(0) segments, and the second of which transfers data between the other three domain(0) segments. Data is transferred between the two domain(1) segments over a domain(2) segment having two domain routing elements, each connected with a corresponding one of the domain(1) segments.\nThe system\"\"s memory elements each include a directory element that maintains a list of descriptors reflecting the identity and state of each datum stored in the corresponding memory. One portion of each descriptor is derived from the associated datum\"\"s system address, while another portion represents an access state governing the manner in which the local central processing unit may utilize the datum. This access state may include any one of an xe2x80x9cownershipxe2x80x9d state, a read-only state, and an invalid state. The first of these states is associated with data which can be modified by the local central processing unit, i.e., that unit included within the cell in which the datum is stored. The read-only state is associated with data which may be read, but not modified, by the local central processing unit. The invalid state is associated with invalid data copies.\nThe domain routing elements themselves maintain directories listing all descriptors stored in their descendant domain(0) segments. Thus, in the above example, the routing elements of first domain(1) segments maintain directories reflecting the combined content of the cells of their respective domain(0) segment. Moreover, the single routing element of the domain(2) segment maintains a directory listing all descriptors retained in all of the system\"\"s processing cells.\nData access requests generated by a processor are handled by the local memory element whenever possible. More particularly, a controller coupled with each memory monitors the cell\"\"s internal bus and responds to local processor requests by comparing the request with descriptors listed in the corresponding directory. If found, matching data is transmitted back along the internal bus to the requesting processor.\nData requests that cannot be resolved locally are passed from the processing cell to the memory management system. The management element selectively routes those unresolved data requests to the other processing cells. This routing is accomplished by comparing requested descriptors with directory entries of the domain routing units. Control elements associated with each of those other cells, in turn, interrogate their own associated directories to find the requested data. Data satisfying a pending request is routed along the domain segment hierarchy from the remote cell to the requesting cell.\nData movement between processing cells is governed by a protocol involving comparative evaluation of each access request with the access state associated with the requested item. The memory management system responds to a request for exclusive ownership of a datum by moving that datum to the memory element of the requesting cell. Concurrently, the memory management element allocates physical storage space for the requested item within the requesting cell\"\"s data storage area. The management element also invalidates the descriptor associated with the requested item within the data store of the remote cell, thereby effecting subsequent deallocation of the physical storage space which had retained the requested item prior to its transfer to the requesting cell.\nWhile the aforementioned operations result in exclusive storage of the requested datum within the requesting cell, other cells may subsequently gain concurrent access to that datum, for example, on a read-only basis. Particularly, the memory management system responds to a request by a first cell for read-only access to datum exclusively owned by a second cell by transmitting a copy of that datum to the first cell while simultaneously designating the original copy of that data, stored in the second cell, as xe2x80x9cnonexclusively owned.xe2x80x9d\nThe system permits an owning cell to disable the copying of its data by providing a further ownership state referred to as the xe2x80x9catomicxe2x80x9d state. The memory management system responds to requests for data in that state by transmitting a wait, or xe2x80x9ctransient,xe2x80x9d signal to requestors and by broadcasting the requested data over the hierarchy once atomic ownership is relinquished.\nA system of the type described above provides improved multiprocessing capability with reduced bus and memory contention. The dynamic allocation of exclusive data copies to processors requiring exclusive access, as well as the sharing of data copies required concurrently by multiple processors reduces bus traffic and data access delays. Utilization of a hardware-enforced access protocol further reduces bus and memory contention, while simultaneously decreasing software overhead required to maintain data coherency. The interconnection of information transfer domain segments permits localization of data access, transfer and update requests.\nThese and other aspects of the invention are evident in the description which follows."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital electronic camera and the interfacing of the camera to an external processing device which monitors, receives images and/or audio, and/or controls the camera through an input/output interface. The invention is more particularly related to a digital electronic camera which automatically detects a connection to the external processing device.\n2. Discussion of the Background\nConventional cameras which use light sensitive film have been known for a long time and the manner of generating images is well-known; light sensitive photographic film is exposed to light which passes through a lens and a shutter. The film is then taken to a developing center and processed into photographic prints.\nDigital electronic cameras which electronically capture images or images along with audio are also known. However, this field of technology is relatively new and there is not a universal standard for exporting or producing images from digital cameras. When designing the camera corresponding to the present invention, the inventors have discovered various problems in communicating information out of the camera. These problems include the connection of the camera to an external device such as another camera or a general purpose computer. The present inventors have sought to overcome problems pertaining to the actual connection of the camera to the external device, the conversion of signal levels from a level used by the processor within the camera to a level compatible with a communication interface, the detection of a connection of the camera to the external device, and the manner of assembling the information to be transmitted."}
{"text": "High-Efficiency Video Coding (HEVC) is a new international video coding standard that is being developed by the Joint Collaborative Team on Video Coding (JCT-VC). HEVC is based on the hybrid block-based motion-compensated DCT-like transform coding architecture. The basic unit for compression, termed Coding Unit (CU), is a 2N×2N square block, and each CU can be recursively split into four smaller CUs until a predefined minimum size is reached. Each CU contains one or multiple Prediction Units (PUs). The PU sizes can be 2N×2N, 2N×N, 2N×nU, 2N×nD, N×2N, nL×2N, nR×2N, or N×N, where 2N×N, 2N×nU, 2N×nD and N×2N, nL×2N, nR×2N correspond to horizontal and vertical partition of a 2N×2N PU with symmetric or asymmetric PU size division respectively.\nTo further increase the coding efficiency of motion vector coding in HEVC, the motion vector competition (MVC) based scheme is applied to select one motion vector predictor (MVP) among a given MVP candidate set which includes spatial and temporal MVPs. There are three inter-prediction modes including Inter, Skip, and Merge in the HEVC test model version 3.0 (HM-3.0). The Inter mode performs motion-compensated prediction with transmitted Motion Vector Differences (MVDs) that can be used together with MVPs for deriving motion vectors (MVs). The Skip and Merge modes utilize motion inference methods (MV=MVP+MVD where MVD is zero) to obtain the motion information from spatial neighboring blocks (spatial candidates) or temporal blocks (temporal candidates) located in a co-located picture. The co-located picture is the first reference picture in list 0 or list 1, which is signaled in the slice header.\nWhen a PU is coded in either Skip or Merge mode, no motion information is transmitted except for the index of the selected candidate. In the case of a Skip PU, the residual signal is also omitted. For the Inter mode in HM-3.0, the Advanced Motion Vector Prediction (AMVP) scheme is used to select a motion vector predictor among an AMVP candidate set including two spatial MVPs and one temporal MVP. In this disclosure MVP may refer to motion vector predictor or motion vector prediction. As for the Merge and Skip mode in HM-3.0, the Merge scheme is used to select a motion vector predictor among a Merge candidate set containing four spatial MVPs and one temporal MVP.\nFor the Inter mode, the reference picture index is explicitly transmitted to the decoder. The MVP is then selected among the candidate set for a given reference picture index. FIG. 1 illustrates the MVP candidate set for the Inter mode according to HM-3.0, where the MVP candidate set includes two spatial MVPs and one temporal MVP:                1. Left predictor (the first available MV from A0 and A1),        2. Top predictor (the first available MV from B0, B1, and Bn+1), and        3. Temporal predictor (the first available MV from TBR and TCTR).        \nA temporal predictor is derived from a block (TBR or TCTR) in a co-located picture, where the co-located picture is the first reference picture in list 0 or list 1. The block associated with the temporal MVP may have two MVs: one MV from list 0 and one MV from list 1. The temporal MVP is derived from the MV from list 0 or list 1 according to the following rule:                1. The MV that crosses the current picture is chosen first, and        2. If both MVs cross the current picture or both do not cross, the MV with the same reference list as the current list will be chosen.        \nIn HM-3.0, if a particular block is encoded in the Merge mode, an MVP index is signaled to indicate which MVP among the MVP candidate set is used for this block to be merged. To follow the essence of motion information sharing, each merged PU reuses the MV, prediction direction, and reference picture index of the selected candidate. It is noted that if the selected MVP is a temporal MVP, the reference picture index is always set to the first reference picture. FIG. 2 illustrates the MVP candidate set for the Merge mode according to HM-3.0, where the MVP candidate set includes four spatial MVPs and one temporal MVP:                1. Left predictor (Am),        2. Top predictor (Bn),        3. Temporal predictor (the first available MV from TBR or TCTR),        4. Above-right predictor (B0), and        5. Below-left predictor (A0).        \nIn HM-3.0, a process is utilized in both Inter and Merge modes to avoid an empty candidate set. The process adds a candidate with a zero MV to the candidate set when no candidate can be inferred in the Inter, Skip or Merge mode.\nBased on the rate-distortion optimization (RDO) decision, the encoder selects one final MVP for Inter, Skip, or Merge modes from the given MVP list and transmits the index of the selected MVP to the decoder after removing redundant MVPs in the list. However, because the temporal MVP is included in the MVP list, any transmission error may cause parsing errors at the decoder side and the error may propagate. When an MV of a previous picture is decoded incorrectly, a mismatch between the MVP list at the encoder side and the MVP list at the decoder side may occur. Therefore, subsequent MV decoding may also be impacted and the condition may persist for multiple subsequent pictures.\nIn HM-4.0, in order to solve the parsing problem related to Merge/AMVP in HM-3.0, fixed MVP list size is used to decouple MVP list construction and MVP index parsing. Furthermore, in order to compensate the coding performance loss caused by the fixed MVP list size, additional MVPs are assigned to the empty positions in the MVP list. In this process, Merge index is coded using truncated unary codes of fixed length equal to 5 or less, and AMVP index is coded using fixed length equal to 2 or less.\nAnother change in HM-4.0 is the unification of MVP positions. Both Merge and Skip use the same positions shown in FIG. 3. For Merge mode in HM-4.0, up to four spatial MVPs are derived from A0, A1, B0, and B1, and one temporal MVP is derived from TBR or TCTR. For the temporal MVP, TBR is used first. If TBR is not available, TCTR is used instead. If any of the four spatial MVPs is not available, the block position B2 is then used to derive MVP as a replacement. After the derivation process of the four spatial MVPs and one temporal MVP, the process of removing redundant MVPs is applied. If the number of available MVPs is smaller than five after redundant MVP removal, three types of additional MVPs are derived and are added to the MVP list.\nIn the derivation for the spatial and temporal MVPs, the MVP can be derived with the MV pointing to the same reference picture as the target reference picture. Alternatively, the MVP can be derived from a candidate MV pointing to a different reference picture. FIG. 4 illustrates an example of deriving a spatial MVP based on various types of motion vectors associated with spatial neighboring candidate blocks, where the candidate blocks comprises spatial neighboring blocks A0, A1, B0, B1 and B2, and temporal co-located blocks TBR or TCTR. The circled numbers refer to the search order for determining an MVP from respective candidates. The highest priority of search corresponds to an MV pointing to the target reference picture within the given reference list. The next highest priority of search corresponds to an MV pointing to the target reference picture within the other reference list. The third and fourth search priorities correspond to other reference picture within the given and other reference lists respectively. In the particular example of FIG. 4, the availability of motion vectors 1 and 2 is checked together, and the availability of motion vectors 3 and 4 is checked together. The availability of motion vectors 1 and 2 is checked from candidate blocks A0 through A1 and then from B0 through B2. If none of the MVs exists, the search checks the availability of motion vector 3 and 4 through all the blocks. When the MVP is derived from an MV pointing to a different reference picture or the MV is for a co-located picture, the MV may have to be scaled to take into consideration of different picture distances. The exemplary search patterns for MVP derivation as shown in FIG. 4 shall not be construed as limitations to the present invention as described in this application. For example, the availability of motion vectors 1 through 4 for each block can be checked together. In another example, motion vector 1 can be checked first in the order from A0 to A1 and then from B0 to B2. If none of the MVs exists, the search moves to check the availability of motion vector 2 from A0 to A1 and then from B0 to B2. The process will continue for motion vector 3 and motion vector 4 if the spatial MVP has not been derived.\nIn the derivation process for the spatial and temporal MVPs, the division operation is required to scale the motion vector. The scaling factor is calculated based on the picture distance ratio. For example, the MVP may be derived based on the MV of a co-located block. The picture distance scaling factor, DistScaleFactor is computed according to:\n                              DistScaleFactor          =                                                    POC                curr                            -                              POC                ref                                                                    POC                temp                            -                              POC                                  temp                  ⁢                  _                  ⁢                  ref                                                                    ,                            (        1        )            where POCcurr and POCref represent the picture order counts (POCs) of the current picture and the POC of the target reference picture respectively, and POCtemp and POCtemp_ref represent the POC of the co-located picture and the POC of the reference picture pointed to by the MV of the co-located block respectively. While the MV of a co-located block is used to illustrate the derivation of picture distance scaling factor, the MV of a spatial neighboring block may also be used to derive the MVP and the corresponding derivation of picture distance scaling factor can be shown similarly.\nIn an implementation according to HM-4.0, the POC distance between current picture and the target reference picture and the POC distance between the co-located picture and the reference picture pointed to by the MV of the co-located block are first constrained within a given range, i.e.:DiffPOCcurr=clip(−128,127,POCcurr−POCref),DiffPOCtemp=clip(−128,127,POCtemp−POCtemp_ref).\nThe scaling factor is then calculated according to the following equations:\n                                              ⁢                              X            =                                                            2                  14                                +                                                                                              DiffPOC                      temp                                        2                                                                                                DiffPOC                temp                                              ,          and                                    (        2        )                                DistScaleFactor        =                  clip          ⁡                      (                                          -                1024                            ,              1023              ,                                                (                                                                                    DiffPOC                        curr                                            ×                      X                                        +                    32                                    )                                >>                6                                      )                                              (        3        )            \nThe scaled MVP is derived by multiplying the MV by the distance scaling factor, i.e.:ScaledMV=sign(DistScaleFactor×MV)×((|DistScaleFactor×MV|+127)>>8)  (4)\nThe picture scaling factor in the form of equation (1) will require division operation, which is more complicated for hardware implement or consumes more CPU time in software based implementation. The essence of the implementation based on HM4.0 is to pre-multiply the distance ratio by a multiplication factor (214 in this equation (2)) so that the distance ratio becomes an integer. In equation (2), an offset term, |DiffPOCtemp/2|is added to 214 to take care of data conversion with rounding. Similarly, offsets 32 and 127 are added in equations (2) and (3) for data conversion. The multiplication factor can be compensated by simple right-shifting operation. Therefore, the implementation associated with equations (1) through (4) is a preferred approach since the computation of the scaled MV, ScaledMV does not require any division operation.\nIn HM-4.0, the scaling factor (DistScaleFactor) is clipped to the range [−1024, 1023], as shown in equation (3). The scaling factor will be right shifted by eight bits as shown in equation (4), which implies that the effective scaling range is limited to [−4, 4). A reference picture selection method for a low-delay coding system is disclosed by Lii et al. to improve coding efficiency (“Encoding optimization to improve coding efficiency for low delay cases”, by Li et al., Joint Collaborative Team on Video Coding (JCT-VC) of ITU-T SG16 WP3 and ISO/IEC JTC1/SC29/WG11 6th Meeting: Torino, IT, 14-22 Jul. 2011, Document: JCTVC-F701). According to the method disclosed by Li, et al., one nearest picture and three high quality pictures (pictures with low QPs) are used as reference pictures for low delay cases. FIG. 5 illustrates an exemplary reference frame configuration for the low-delay system, Picture 10 is the current picture and picture 9 is the nearest picture. Pictures 0, 4 and 8 are the three high quality reference pictures. Block 525 corresponds to a current block and block 515 corresponds to a neighboring block. The neighboring block 515 has an associated MV 510 that is used to derive the MVP for the current MV 520 of the current block 525. The picture distance associated with the current MV is 6 while the picture distance associated with the candidate MV is 1. Therefore, the picture scaling factor for this example is 6, which exceeds the supported scaling factor range. Therefore, the effective scaling range [−4, 4) is not sufficient for some applications.\nAccordingly, it is desirable to develop a scheme to increase the effective scaling ratio for MV scaling. A system incorporating the increased scaling factor may achieve improved performance."}
{"text": "A power supply system that coverts a DC voltage having a given voltage value into a DC voltage having a voltage value different from the given voltage value is known as so-called DC/DC converter. A DC/DC converter is used in various electronic devices, one example of which is a computer. In a computer, for example, a supply voltage supplied to a microprocessor (hereinafter “CPU”) is generated by a DC/DC converter.\nA power supply system like a DC/DC converter has multiple switch elements that cyclically change the direction of a current flowing from one coil to another coil. Each of the multiple switch elements is composed of a transistor, such as metal-oxide-semiconductor field effect transistor (hereinafter “MOSFET”). Various structures of the MOSFET used in the DC/DC converter have been proposed. Patent Documents 1 and 2 disclose various examples of MOSFET structures. In Patent Document 1, a MOSFET structure is shown in, for example, FIG. 1D in which a first gate electrode 10 and a second gate electrode 12 are stacked vertically. In Patent Document 2, another MOSFET structure is shown in, for example, FIG. 4K in which a gate electrode 26 (hereinafter “first gate electrode”) and a gate electrode 30 (hereinafter “second gate electrode”) are stacked vertically."}
{"text": "The present invention relates to electromechanical valves actuators and, more particularly, to electromechanical valve actuators that reduce noise, vibration, and harshness issues associated with the opening of the valve.\nAs engine technology advances and manufacturers strive to increase engine power, improve fuel economy, decrease emissions, and provide more control over engines, manufacturers are developing electromechanical valve actuators to replace cam shafts for opening and closing engine valves. Electromechanical valve actuators allow selective opening and closing of the valves in response to various engine conditions.\nElectromechanical valve actuators generally include two electromagnets formed from a core having an embedded power coil. A spring loaded armature located between the electromagnets is movable between the electromagnets as the power coils are selectively energized to create a magnetic force to attract the armature. The surface of the electromagnets to which the armature is attracted when the power coil of an electromagnet is energized is generally referred to as a pole face. The armature abuts to the valve so that as the armature moves between pole faces in pole-face-to-pole-face operation, the valve is opened and closed.\nWhen the valve is in a closed position, the armature plate is generally held against or near the armature electromagnet and a gap is created between the armature stem and the valve stem. This gap is commonly referred to as a lash gap. The lash gap allows thermal movement of metal parts during engine operation and is necessary to ensure that under all conditions the valve is fully closed, while the armature is seated on or near the armature electromagnet pole face. One problem with traditional electromechanical valves is that noise, vibration, and harshness issues occur when the armature stem contacts the valve stem as the valve is opened. More specifically, as the armature plate is released from the armature electromagnet, the armature spring pushes the armature assembly, specifically the armature stem, toward the valve stem so that the armature stem, contacts the valve stem, typically a lash cap on the valve stem, at a high velocity to move the valve from the closed position to the open position. The impact between the valve stem and armature stem may cause noise, vibration, and harshness issues. These noise, vibration, and harshness issues may be amplified due to the excitement of other components within the electromechanical valve actuator, such as the armature plate which may act as a radiator to amplify the noise.\nWhile some manufacturers have attempted to alleviate these problems by combining traditional hydraulic lifters with electromechanical valve actuators, problems may occur from this combination. Hydraulic lifters typically increase the friction experienced by electromechanical valve actuators and add mass to the moving parts. Any increase in mass or friction requires additional power consumption by the electromechanical valve actuator to move the armature and valve and such additional power consumption is particularly acute in opening the valve during the exhaust cycle. Additional power consumption raises additional issues such as undesirable excess heating of the electromagnet power coils and consumption of additional electrical power from the vehicle's generating system which reduces fuel economy and adds to the cost of the generating and distribution system. Any requirement of additional power puts increased demand on today's already overloaded vehicle electrical systems. Another disadvantage to using traditional hydraulic lifters is that hydraulic lifters are relatively expensive and add to the overall cost of the vehicle.\nTo avoid additional expense, increased friction, increased mass, and increased power consumption associated with hydraulic lifters, some manufacturers have attempted to limit the impact force between the valve stem and armature stem by controlling the current profile supplied to the electromagnets. One method of controlling the current profile includes slowly decreasing the current supplied to the armature electromagnet. For example, the electromechanical valve actuator control system slowly lessens the amount of magnetic attraction of the armature electromagnet. As the force of the armature spring surpasses the combined force of the valve spring and magnetic attraction force of the armature electromagnet, the armature accelerates away from the armature electromagnet. The armature electromagnet continues to apply a magnetic attraction force to the armature plate, thereby slowing the acceleration away from the armature electromagnet. While current profile control reduces noise, vibration, and harshness issues associated with the contact between the armature stem and valve stem during the opening of the valve, current profile control requires additional power consumption as the armature electromagnet works against the force provided by the armature spring. The power consumption of the valve electromagnet is also significantly increased due to the increased magnetic force needed to attract and retain the armature plate against the valve electromagnet. More specifically, because the force applied by the armature spring was inhibited by the armature electromagnet during the initial portion of the open cycle, the valve electromagnet must compensate for this inhibition by increasing the magnetic attraction force of the valve electromagnet. As stated above, to increase magnetic attraction, increased power consumption is required which raises additional problems. Further, the increased power consumption is particularly acute due to the distance between the valve electromagnet and armature plate at which the valve electromagnet must apply the magnetic force especially due to the exponential decrease in magnetic force as the distance increases. Another problem associated with current profile control include the necessity of additional and expensive position sensors and microprocessors to accurately determine and control the position and movement of the armature assembly. Yet another problem with current profile control is that as the armature plate approaches the pole face, the gap between the pole face and armature plate, generally referred to as the air gap, decreases, causing the magnetic force acting on the armature to exponentially increase. This exponential increase is particularly acute due to the additional power being applied by the valve electromagnet to the armature in response to reduced force from the armature spring to ensure attraction and retention of the armature plate as desired. Any increase in magnetic force causes the armature to increase in velocity as it approaches the pole face of the energized electromagnet thereby increasing the force of the impact of the armature. This impact causes noise, vibration, and harshness concerns, also making quiet operation of electromechanical valve actuators challenging to achieve.\nIn view of the foregoing, there is a need for an electromechanical valve actuator having reduced noise, vibration, and harshness issues during opening of the valve, especially during the exhaust cycle, without having increased power consumption."}
{"text": "As televisions and other video displays continue to evolve, the desire for enhanced functionality continues. For example, many computer users now use multiple displays so that different sets of data can be viewed in series or in parallel. These multiple computer displays can typically be organized with a single screen's image being repeated across multiple displays or as an extended desktop wherein the multiple displays are virtually connected wherein the mouse icon will move from one display to another with relative screen locations being determined during setup of the multiple display management software.\nIn addition, many restaurants, open public areas, and other establishments attempt to increase the size of a display by aligning several displays in an attempt to create a larger image. Examples of such technology have been developed by CineMassive® and 9× Media®, Inc. Unfortunately, no previous approach has been successful in aligning multiple displays without either different images being displayed and/or edge line from the displays appearing in the image in a grid-like or other fashion. Such edge lines cause the image to be segmented or non-continuous in nature, which decreases its image quality and its ease of viewing to the user. In view of the foregoing, there exists a need for a solution that solves at least one of the deficiencies in the prior art.\nCurrent LED device pixels have a blocking dam around the edge line of each pixel to protect against color diffusion between pixels. Light diffusion between pixels results in a resolution performance decrease. A modular display with blocking dams around each pixel can bring a visible edge line into the produced image.\nAs LEDs continue to grow in popularity as an efficient technological device, the need for continued advancement grows as well. Along these lines, obtaining white light output from LED is not only needed, but also difficult to achieve. Many approaches in the past have attempted to find new ways to obtain white light. However, many of these approaches perform such processing at the chip level instead of at the wafer level. Such an approach can result in chip waste. Moreover, none of the existing approaches vary phosphor ratios based on an underlying device measurement (such as a wavelength of a light output). For example, U.S. Pat. No. 6,650,044 forms a pedestal on top of a pad. The pedestal is a stud bump that is used for connectivity. This approach is not ideal as it is inefficient, does not provide chip level coating, is much harder to polish, and is easily contaminated. Moreover, in previous approaches such as U.S. Pat. Nos. 7,446,733 and 7,190,328, pads were manufactured so that all LEDs thereon had the same height. In view of the foregoing, there exists a need for a LED and associated method of manufacture that addresses the deficiencies of the related art."}
{"text": "There are several printing methods, and planographic printing method is predominated at present. The planographic printing method is a method in which, employing a planographic printing plate with ink receptive image portions and ink repulsive non-image portions present at substantially the same plane, printing is carried out by supplying ink only to the image portions, and transfer the ink to material to be printed such as paper. That is, planographic printing method employs the difference of ink receptivity between image portions and non-image portions. In order to prepare such a planographic printing plate, a presensitized planographic printing plate (hereinafter referred to also as PS plate) is used. As described in T. Yonezawa, \"PS Ban Gairon\", pp. 18-81, Insatsu Gakkai Shuppanbu (1993), the PS plate herein referred to implies the following:\nThe PS plate implies a presensitized planographic printing plate comprising a hydrophilic aluminum support and provided thereon, a lipophilic light sensitive layer. The planographic printing plate is prepared by remaining the light sensitive layer at image portions and removing the light sensitive layer at non-image portions to expose the hydrophilic aluminum surface capable of forming a dampening water layer and repelling ink, according to lithography. There are a conventional PS plate, from which a planographic printing plate having a dampening water layer as an ink repelling layer is prepared, and a waterless PS plate, from which a waterless planographic printing plate having a silicone rubber layer as an ink repelling layer is prepared.\nIn the conventional PS plate, a grained aluminum plate is ordinarily employed as a support. The surface of the aluminum plate is required to have a water holding property and an excellent adhesion to a lipophilic light sensitive layer, such that the light sensitive layer is not removed from the surface during printing. Therefore, the surface of the aluminum plate is ordinarily grained, and optionally subjected to surface treatment such as anodizing treatment to improve both water holding property and adhesion to a lipophilic light sensitive layer. The conventional PS plate is widely employed.\nThe conventional PS plate provides a practical excellent printing plate. However, with conventional printing plates, which employ a grained aluminum plate as the support, it is difficult to control the balance between ink and a dampening water during printing, which is likely to cause printing failure such as stain occurrence and ink emulsification according to printing conditions. Further, such printing plates occasionally result in lowering of small dot reproduction due to halation, which is caused by the surface-roughened aluminum plate.\nThere are two plate making processes in the printing field, one is a direct plate making process of preparing a direct printing plate by directly forming an image necessary to print on a plate, and the other an indirect plate making process of preparing an indirect printing plate by forming an image on an intermediate medium such as a drum, and then transferring the formed image to a plate.\nAs desk top publishing (DTP) systems progress or computer to plate (CPT) plate making process develops, the direct or indirect plate making process has been increasingly used. As one embodiment of the direct plate making processes, there is a method comprising the steps of imagewise transferring a lipophilic transfer layer to the roughened surface of a support to prepare a printing plate, however, the roughened surface occasionally causes problems in that the transferred small dots have fringes around the periphery, which result in print images of poor quality.\nAs one embodiment of the indirect plate making processes, there is a method of preparing an off-set master for a printing plate. The off-set master has an image receiving layer comprised of zinc oxide and a binder on a paper support or a plastic film support. Toner images are formed on the image receiving layer employing a copy machine or a laser printer to prepare an off-set printing plate. The off-set printing plate has lipophilic, toner image portions and non-image portions without toner images, which are made to be hydrophilic by an etching solution containing, for example, potassium cyanide.\nThe off-set master requires etching before printing, whereby non-image portions of the master are made to be hydrophilic. However, the etching treatment still has a problem in that it is insufficient to make the non-image portions hydrophilic, and the master is likely to produce stains at the non-image portions.\nAs is described above, a support for a planographic printing plate having satisfactory properties has not yet been obtained."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a semiconductor integrated circuit. More particularly, the present invention relates to an apparatus and method for preventing a system error caused by circuit synchronization and control disable circumstances upon the generation of a glitch when switching clocks.\n2. Description of the Related Art\nSemiconductor integrated circuits use clocks for the synchronization and control of logic circuits, flip-flop circuits and the like. In recent years, semiconductor integrated circuits provide a diversity of functions. The semiconductor integrated circuits integrating a diversity of functions include certain functions that require more clock sources than other functions, thus requiring a process of switching between many clocks.\nA conventional clock switching technology performs switching without addressing a glitch of a clock caused by a select signal of a multiplexer. However, if a glitch is generated upon the switching of clocks, such a glitch has adverse effects upon a system, often making the system unstable.\nFIGS. 1A and 1B are diagrams illustrating the generation of a glitch in a conventional clock switching circuit.\nIn FIG. 1A, a clock switching circuit outputs one of two clock signals, i.e., a first clock signal 100, or a second clock signal 102, by using a select signal (SEL) 104 for clock selection.\nIn FIG. 1B, a clock switching method can output either a first clock 110, or a second clock 120, depending on a value of select signal 130. When select signal 130 is low, clock signal 110 (clk 1) is output. When select signal 130 is high, clock signal 120 (clk 2) is output.\nHowever, upon switching between the clocks 110 and 120, a glitch 140 is generated by a timing difference between the clocks 110 and 120.\nIn other words, upon switching between the clocks 110 and 120, a glitch is typically generated and may give rise to circuit synchronization problems and control disable problems, thus having a very bad impact upon a system.\nMethods for preventing the generation of a glitch have been proposed to solve the glitch problem. However, these methods are mostly technologies for switching between two clocks, or switching to avoid generating a glitch considering a clock switching time or a relationship between the clocks.\nFor example, there has been previously proposed a method for removing a glitch by counting a predetermined value during a first clock and a second clock change interval while fixing a clock output CLOCK_OUT. This method is suitable when there is no change of a frequency difference between two clocks and switching between two clocks. However, the aforementioned method has problems when there is a large frequency difference between two clocks, or when there is switching between three or more clocks.\nIn addition to the above method, there is also a method that operates upon detecting a change of a select signal, and in response, gating a first clock and synchronizing a select signal of a multiplexer to a second clock that is a second clock, thus enabling the clock switching and removing a glitch. The above method also has problems when the number of clocks is three or more."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and system for distributing the same data to a plurality of nodes constructing a PBX (Private Branch Exchange) network.\n2. Description of the Related Art\nHitherto, when the same local data is desired to be used by nodes in a PBX network constructed by a plurality of nodes, it is troublesome since a maintenance terminal has to access each of the nodes and enter the same data to each node. As the number of nodes increases, it becomes more troublesome.\nIt is an object of the invention to provide a data distributing method and system which can efficiently distribute local data supplied from a single maintenance terminal to all of nodes and certainly synchronize data among the nodes in a PBX network comprised of a plurality of nodes.\nAccording to a data distributing method of the invention, one of a plurality of nodes constructing the PBX network is set as a master node and the other nodes are set as slave nodes, and local data supplied from a data input unit to the master node is retained in a data retaining unit in the master node. In the master node, the local data is formed in a transfer format by adding a checksum and a destination node designated to each slave node to the local data, and the local data in the transfer format is transferred to the slave node. In each slave node, the local data is retained into a data retaining unit in the slave node after checking whether the designated destination node is the slave node itself or not and checking the checksum.\nThe data transfer from the master node to each of the slave nodes is performed by broadcast transfer through a transmission line which can perform broadcast transfer.\nEach of the slave nodes retains the local data in its data retaining unit, after that, returns reply data formed in a reply format in which the node number of the slave node is set as a transmission source and the node number of the master node is set as a destination to the master node, and the master node receives the reply data and transfers it to the data input unit.\nThe master node makes the data length of the local data included in the transfer format and transfers the local data in the transfer format.\nThe master node divides the local data every predetermined length into blocks, makes a serial number of each block included in the transfer format, and transfers the data in the transfer format.\nThe same data as the local data which is transferred to the slave node and stored in the data retaining unit in the slave node is similarly transferred for data collation from the master node, and each of the slave nodes collates the local data newly transferred with the local data which has been already retained in its data retaining unit and returns a collation result in a reply format in which the node number of the slave node is set as a transmission source and the node number of the master node is set as a destination to the master node.\nAccording to a data distributing system of the invention, one of a plurality of nodes constructing a PBX network is set as a master node and the other nodes are set as slave nodes, the master node comprises a master data retaining unit, data writing means for writing local data supplied from a data input unit into the master data retaining unit, data editing means for editing the supplied local data together with a checksum of the local data and a destination node designated to each of the slave nodes into a transfer format, and data transfer means for broadcast transferring the edited data through a transmission line which can perform broadcast transfer to the slave node. Each of the slave nodes comprises a slave data retaining unit, data receiving means for checking whether data transferred from the master node is addressed to the slave node itself or not and receiving the data, and data writing means for writing the local data into the slave data retaining unit when the checksum of the received transfer data is checked and matched.\nEach of the slave nodes comprises data return means for forming the data in a reply format in which the node number of the slave node itself is set as a transmission source and the node number of the master node is set as a destination after the local data is retained in the slave data retaining unit of the slave node and returning reply data in the format to the master node, and the master node has data receiving means for receiving the data returned from the slave node and transferring the data to the data input unit.\nThe data editing means in the master node makes the data length of the local data included in the transfer format.\nThe data editing means in the master node divides the local data into blocks on a predetermined length unit basis and makes a serial number of each block included in the transfer format.\nThe data transfer means in the master node similarly transfers the local data for data collation, which is the same as the local data transferred to the slave node and stored in the slave data retaining unit in the slave node. Each of the slave nodes is further provided with a data buffer unit for storing the local data newly transferred and data collating means for collating the local data stored in the data buffer unit with the local data which has been already retained in the data retaining unit, and the data returning means in each of the slave nodes returns a collation result in a reply format in which the node number of the slave node itself is set as a transmission source and the node number of the master node is set as a destination to the master node."}
{"text": "The present invention relates to telecommunication systems and more particularly to digital hierarchy systems (SDH). In such systems data is switched within an individual station to the appropriate output port. Also alternate routes are provided between stations by which data can be transmitted from an originating stations port to a receiver stations port so that failure of one route, e.g. due to accidental damage, will not prevent the data from reaching the receiver stations port via the other route.\nThe data is sent in the form of two packages or so-called virtual containers which essentially are made up of two sections, one being the actual data which is transmitted, i.e. the soiled payload, and the other being data which is concerned with the integrity of the data, referred to hereinafter as overhead or ancillary data.\nOne of the requirements for a system of the kind outlined above is that there must be some means for checking whether the package is correctly switched (or routed) at each individual station along the overall path of the data. Essentially this is achieved by use of the ancillary or overhead data area at each station in order to include an identity, and other monitoring information, to ensure correct switching.\nIn our co-pending application no. 9301575 we disclose an arrangement where the package is identified in terms of the channel number of the data, and the receiving stations card number. However there is a potential problem with such an arrangement in situations where the package travels by the alternative route referred to earlier. In the event of failure of the first route the package will acquire a new channel and card number. This means that the destination port of the station will not recognise the package because it will be looking for a package which has an identification which is characteristic of the first route. This problem is a direct result of the use of the alternate route protection identified earlier.\nThe present invention is concerned with overcoming this problem."}
{"text": "Resistance variable memory elements, which include chalcogenide-based programmable conductor elements, have been investigated for suitability as semi-volatile and non-volatile random access memory devices. A typical such device is disclosed, for example, in U.S. Pat. No. 6,849,868 to Campbell, which is incorporated by reference.\nIn a typical chalcogenide-based programmable conductor memory device, a conductive material, such as silver, is incorporated into a chalcogenide glass. The resistance of the chalcogenide glass can be programmed to stable higher resistance and lower resistance states. An unprogrammed chalcogenide-based programmable conductor memory device normally in a higher resistance state. A write operation programs the chalcogenide-based programmable conductor memory device to a lower resistance state by applying a voltage potential across the chalcogenide glass. The chalcogenide-based programmable conductor memory device may then be read by applying a voltage pulse of a lesser magnitude than required to program it; the resistance across the memory device is then sensed as higher or lower to define the ON and OFF states.\nThe programmed lower resistance state of a chalcogenide-based programmable conductor memory device can remain intact for an indefinite period, typically ranging from hours to weeks, after the voltage potentials are removed. The chalcogenide-based programmable conductor memory device can be returned to its higher resistance state by applying a reverse voltage potential of about the same order of magnitude as used to write the device to the lower resistance state. Again, the higher resistance state is maintained in a semi- or non-volatile manner once the voltage potential is removed. In this way, such a device can function as a variable resistance memory having at least two resistance states, which can define two respective logic states, i.e., at least a bit of data.\nOne exemplary chalcogenide-based programmable conductor memory device uses a germanium selenide (i.e., GexSe100−x) chalcogenide glass as a backbone. The germanium selenide glass has, in the prior art, incorporated silver (Ag) and silver selenide (Ag2Se).\nPrevious work by the inventor, Kristy A. Campbell, has been directed to chalcogenide-based programmable conductor memory devices incorporating a silver-chalcogenide material as a layer of silver selenide (e.g., Ag2Se) or silver sulfide (e.g., Ag2S) in combination with a silver-metal layer and a chalcogenide glass layer. The silver-chalcogenide materials are suitable for assisting in the formation of a conducting channel through the chalcogenide glass layer for silver ions to move into to form a conductive pathway.\nTin (Sn) has a reduced thermal mobility in Ge2Se100−x compared to silver and the tin-chalcogenides are less toxic than the silver-chalcogenides, therefore tin-chalcogenides (e.g., SnSe) have also been found to be useful in chalcogenide-based programmable conductor memory devices to replace silver selenide. However, sputtering of tin selenide to form such devices has proven difficult due to the increased density of the sputtered layers. This increased density (e.g., ˜6 g/cm3 sputtered compared to ˜3 g/cm3 evaporated) can prevent the motion of silver ions into the chalcogenide glass, thereby preventing the memory device from functioning. Therefore, evaporative deposition techniques have been used to deposit such material, which is generally a less efficient, more costly, slower, and less controlled technique for deposition. However, evaporation deposition of tin selenide and silver also incorporates some oxygen into the resulting layer, which provides for the lower density and allows for more mobility of silver ions."}
{"text": "The goal of plant breeding is to combine, in a single variety or hybrid, various desirable traits. For field crops, these traits may include resistance to diseases and insects, resistance to heat and drought, reducing the time to crop maturity, greater yield, and better agronomic quality. With mechanical harvesting of many crops, uniformity of plant characteristics such as germination, stand establishment, growth rate, maturity, plant height and ear height, is important. Traditional plant breeding is an important tool in developing new and improved commercial crops."}
{"text": "1. Technical Field:\nThe present invention relates in general to an improved means for monitoring the physical condition of a data recording device, and in particular to an improved means for identifying damage to the components of a data recording device. Still more particularly, the present invention relates to an improved apparatus for detecting fractures in a recording head or tape guide in a data tape drive.\n2. Description of the Prior Art:\nData recording devices, such as data tape drives, record information to or read information from media in a storage device, such as the data tape of a tape cartridge. As shown in FIG. 1, data tape drives utilize one or more transducer heads 1 having air bleed grooves or slots 3 that define side walls 5. The air bleed slots 3 are provided for optimizing the performance of the transducer head 1. For example, U.S. Pat. No. 5,636,085 discloses a magnetic read/write assembly for a flexible disk drive having rails with air bleed slots that enable a low flying height for the media relative to the electromagnetic elements.\nReferring again to FIG. 1, air bleed slots 3 are typically formed in the transducer head 1 in a direction that is perpendicular to the direction of travel of the tape (indicated by arrows 7). However, the slots also may be formed or cut in directions parallel to the direction of travel of the magnetic tape. The air bleed slots allow the tape to ride closer to the read and write elements in the head by bleeding off the boundary layer air which sticks to the moving tape. Although tape tension is used to wrap the tape over the head, high density recording requires the tape to hydrodynamically fly extremely close to the recording elements.\nTransducer heads and the tape guides adjacent to them are typically formed from sensitive materials that are somewhat brittle in nature, and are subject to brittle fracture damage due to vibration, shock, and incidental contact with the media. When air bleed slots or other mechanical formations are formed in transducer heads and tape guides, areas of stress concentration occur and may increase the likelihood of damage to these components. The side walls of the transducer head are particularly vulnerable to stress concentrations and fractures. When a transducer head or tape guide is damaged, the media or data tape may be exposed to the sharp edges 9 (FIG. 1) of the fracture. Unless the situation is quickly remedied, extensive amounts of the media can be damaged or destroyed, and valuable information can be permanently lost. Thus, an apparatus for detecting damage to the transducer head and tape guide of a data recording device is needed.\nA data tape drive for recording information on the data tape of a tape cartridge has a transducer head with air bleed slots that define side walls, a tape guide, and a device for detecting damage to the transducer head and tape guide. The damage detection device has a conductor loop that is located along the outer edges of the air bleed slots and tape guide for detecting any brittle fracture damage that may be present. The conductor loop is a thin filament of wire that is preferably attached to the ends of the side walls. Alternatively, the filament is mounted along the side walls just below the top surface of the side walls to avoid contact with the moving data tape during operation of the data tape drive. When the transducer head or tape guide experiences a fracture, such as a fracture in one of the side walls, the delicate filament breaks and forms an open circuit. The open circuit in the conductor loop is detected by the data tape drive so that remedial actions can be taken. Without the damage detection device, large quantities of the data tape moving past the transducer head and/or tape guide would be scraped over or gouged by the fracture and permanently damaged, including loss of the information recorded on the data tape.\nThe foregoing and other objects and advantages of the present invention will be apparent to those skilled in the art, in view of the following detailed description of the preferred embodiment of the present invention, taken in conjunction with the appended claims and the accompanying drawings."}
{"text": "In general, heat is an enemy for the IC chip and it is necessary that the internal temperature thereof does not exceed the maximum allowable junction temperature. The electric power consumption per operation area is large in the semiconductor device such as a power transistor or a semiconductor rectifier element. Therefore, the generated heat amount cannot be sufficiently released with only the heat amount released from a case (package) and a lead of the semiconductor device. It is feared that the internal temperature of the device may be raised to cause thermal destruction.\nThis phenomenon also occurs in the same manner in the IC chip which carries a CPU. The amount of heat generation is increased during the operation in proportion to the improvement in clock frequency. It is an important matter to make the thermal design in consideration of the heat release.\nIn the thermal design for preventing the thermal destruction or the like, element design or package design is made on condition that a heat sink having a large heat release area is secured to a case (package) of the IC chip.\nIn general, a metal material such as copper and aluminum, which has a good thermal conductivity, is used as a material for the heat sink.\nRecently, the IC chip such as CPU and memory is in a trend that the IC chip itself has a large size in accordance 10, with the high degree of integration of the element and the enlargement of the element-forming area, while it is intended to drive the IC chip at low electric power for the purpose of low electric power consumption. When the IC chip has such a large size, it is feared that the stress caused by the difference in thermal expansion between the semiconductor substrate (silicon substrate or GaAs substrate) and the heat sink is increased, and that the peeling-off phenomenon and the mechanical destruction occur in the IC chip.\nIn order to avoid such an inconvenience, for example, it may be pointed out that the low electric power driving of the IC chip should be realized, and the heat sink material should be improved. The low electric power driving of the IC chip is realized in the level of not more than 3.3 V at present and the TTL level (5 V) which has been hitherto used as the power source voltage becomes obsolete.\nAs for the constitutive material for the heat sink, it is insufficient to consider only the thermal conductivity. It is necessary to select a material which has a coefficient of thermal expansion approximately identical with those of silicon and GaAs, which are used as the semiconductor substrate, while having a high thermal conductivity at the same time.\nA variety of reports have been made in relation to the improvement of the heat sink material, including, for example, a case in which aluminum nitride (AlN) is used and a case Cu (copper)-W (tungsten) is used. AlN is excellent in balance between the thermal conductivity and the thermal expansion. Especially, the coefficient of thermal expansion of AlN is approximately coincident with the coefficient of thermal expansion of Si. Therefore, AlN is preferred as a heat sink material for a semiconductor device in which a silicon substrate is used as the semiconductor substrate.\nCu—W is a composite material having both of the low thermal expansion of W and the high thermal conductivity of Cu. Further, Cu—W is mechanically machined with ease. Therefore, Cu—W is preferred as a constitutive material for a heat sink having a complicated shape.\nOther instances have been suggested, wherein metal Cu is contained in a ratio of 20 to 40% by volume in a ceramic base material containing a major component of SiC (conventional technique 1, see Japanese Laid-Open Patent Publication No. 8-279569), and wherein a powder-sintered porous member of an inorganic substance is infiltrated with Cu by 5 to 30% by weight (conventional technique 2, see Japanese Laid-Open Patent Publication No. 59-228742).\nThe heat sink material concerning the conventional technique 1 is produced in the powder formation in which a green compact of SiC and metal Cu is formed to produce a heat sink. Therefore, the coefficient of thermal expansion and the coefficient of thermal conductivity represent only theoretical values. It is impossible to obtain the balance between the coefficient of thermal expansion and the coefficient of thermal conductivity required for actual electronic parts etc.\nThe conventional technique 2 uses a low ratio of Cu with which the powder-sintered porous member composed of the inorganic substance is infiltrated. It is feared that a limit may arise to enhance the thermal conductivity thereby.\nOn the other hand, a composite material, which is obtained by combining carbon and metal, has been developed and practically used. However, such a composite material is used, for example, as an electrode for the discharge machining when the metal is Cu. When the metal is Pb, such a composite material is used, for example, as a bearing material. No case is known for the application as a heat sink material.\nThat is, in the present circumstances, the coefficient of thermal conductivity is at most 140 W/mK for the composite material obtained by combining carbon and metal, which cannot satisfy the value of not less than 160 W/mK required for the heat sink material for the IC chip."}
{"text": "Five unique human hepatitis viruses have been identified (1-5). The hepatitis A virus and hepatitis E virus are enterically transmitted RNA viruses that do not cause chronic liver disease. In contrast, the hepatitis B virus, hepatitis C virus and hepatitis D virus (HBV, HCV and HDV, respectively) are parenterally transmitted and cause chronic infection. They are dangerous contaminants of the blood supply. Recently tests have become readily available for testing for HBV in blood, allowing for the screening for this pathogen and the elimination of infected samples from the blood supply (6).\nConcomitant with the availability of the HBV test came an increase in the proportion of cases of post-transfusion hepatitis due to non-A, non-B (NANB) agents. Until recently, there was no test available for the detection of the NANB agents. The principal NANB agent, HCV, was recently identified by molecular cloning of segments of the HCV genome (3). HCV is an RNA virus related to human flaviviruses and animal pestiviruses (7,8).\nProspective studies of selected counties in the United States by the Centers for Disease Control (CDC) indicate that approximately 170,000 new cases of NANB/HCV infection occur yearly (9). At least 50% of these infections appear to progress to chronic liver disease. Severe sequelae include the development of decompensated cirrhosis necessitating liver transplantation, and development of hepatocellular carcinoma (10,11).\nThe positive-stranded RNA genome of the HCV contains approximately 10,000 nucleotides. The HCV genome acts as a long open reading frame (ORF) capable of encoding a 3,010 amino acid polyprotein precursor from which individual viral proteins, both structural and nonstructural, are produced (7,12-14). There are at least 324 nucleotides at the 5'-end of the ORF which have not yet been shown to encode for protein. Thus, this sequence is referred to as the 5'-non-coding region (7,12-17). Several research groups have reported the nucleotide sequence of either the whole HCV genome or specific subgenomic regions (7,12-22). Comparison of these sequences demonstrates variations in the structural and nonstructural regions (ranging from 9-26%) among different HCV strains. In contrast, the sequences of the 5'-non-coding region appear to have a homology of approximately 99% among different strains (16,17). The 5'-non-coding region also has substantial homology (45-49%) with the equivalent region of animal pestiviruses (7).\nTwo major techniques are currently used to detect HCV infection. The first technique detects antibody produced in response to HCV infection (anti-HCV) (23-28). Since multiple weeks are required for infected patients to develop detectable IgG antibody against HCV antigens, this test is useless in the detection of acute HCV infection. Moreover, studies indicate that antibody testing is associated with both false positive and false negative results (29).\nThese shortcomings in the original assays have spurred development of newer supplemental antibody tests for the diagnosis of HCV infection (30). Preliminary results with supplemental assays indicate a decrease in the frequency of false-positive and negative results. However, false-positive and -negative results still occur and supplemental tests remain unsuitable for detection of acute infection (31).\nThe second technique, detection of HCV RNA by an RNA polymerase chain reaction (PCR), has been limited to research use. The HCV PCR evaluates infection by detecting HCV RNA in blood or tissue extracts through reverse transcription and cDNA amplification (7,32-41). HCV PCR represents a sensitive, direct technique but requires meticulous care (7) to prevent false positive and negative results. The HCV PCR technique, in contrast to antibody tests, can detect circulating HCV RNA during acute infection.\nThe original HCV PCR tests used primers specific for sequences in the non-structural region of the HCV genome (32-36). Subsequently, HCV PCR has been performed using several primers for the 5'-non-coding region in the genome (37,39). In our laboratory we have established HCV PCR for both the nonstructural and 5'-non- coding regions. Our comparative results indicate that the HCV PCR from the 5'-non-coding region is more sensitive in detecting HCV infection (41).\nDespite the success of HCV PCR, the technique has many inherent limitations. First, it is time consuming, expensive and dependent upon meticulous technique. The exquisite sensitivity of PCR makes false positive results due to contamination with exogenous HCV RNA a constant concern (42). Moreover, the variation in both the reverse transcription of HCV RNA to CDNA and the amplification of cDNA make the HCV PCR difficult to quantitate (38,40,42). Recent attempts to overcome these obstacles have resulted in, at best, semi-quantitative assays (38). More importantly, in our experience the efficiency of HCV PCR depends in large part on the specific primers employed. Not only have standards for primers not been developed, but polymerases employed in PCR have different efficiencies. Thus, it will likely be difficult to compare PCR results among different laboratories.\nTo overcome the limitations of current antibody and HCV PCR techniques for detection of HCV infection, it is desirable to develop a test which is highly sensitive, specific, affordable and applicable to the testing of large populations of patients or blood donors. The optimal test would be capable of detecting both acute and chronic infection. Moreover, it would be quantitative to provide information regarding both natural history and the efficacy of current or future antiviral therapies. It would be capable of uniform results. These prerequisites can be fulfilled by a technique to directly detect HCV RNA."}
{"text": "There are many methods for reproducing stereo-pictures (three dimensional pictures) such that the observer has the full spatial impressions of the object. This is done, e.g., by arranging the two pictures recorded from two different directions, as the two eyes of an observer view from the corresponding directions, on a screen and viewing the pictures through glasses of different color for non-color reproduction. If a color image is required, the different pictures are illuminated with linear polarized light of perpendicular orientation and glasses are used with the appropriate polarization such that the eyes see the objects separately in stereo. One difficulty with the foregoing is that the eyes are focused in an unnatural way, i.e. looking in different directions on alternate frames, if the reproduced pictures have to be next to each other.\nThe ideal method requires that the two reproduced pictures of different polarization be located within the same frame. This can be realized, e.g., by holographic production of the color TV picture. This method, however, is basically limited by the coherence length of the radiation such that the stereo effects can be seen only within a certain depth, e.g., 30 centimeters to 1 meter with the current generation of holographic reproduction equipment. This depth, even if extended, is limited by basic principles of physics. A further difficulty is that the granulation of the transmitted picture requires extremely high resolution and very small grain photographic material must be used. The electronic transmission of such pictures also requires an extraordinarily large bandwidth and minimal signal distortion."}
{"text": "1. Field of the Invention\nThis invention relates generally to biohazard surveillance systems and more particularly to an adaptive distributed system for the collection and sampling of hazardous particulates.\n2. Description of the Related Art\nThe challenges we face from biological threat agents are increasing. While microbes continue to evolve and biotechnology becomes more powerful, the inherent hazards to humans, plants, and animals from infectious microorganisms are greatly increased by their intentional use by terrorists. The need for faster and better capabilities for warning, response, and cleanup was painfully evident in the case of a small-scale deployment of a noncontagious, naturally occurring anthrax pathogen. Terrorist use of other biological agents may result in far greater loss of life; agents that might be contagious or perhaps engineered for increased virulence and resistance to medical treatment. As microbes evolve and compete for survival, naturally emerging threats must also be quickly identified and distinguished from suspected terrorism. While the focus on bioterrorism is driven primarily by concerns about attacks on humans, attacks on livestock and/or crops can be just as devastating. A recent outbreak of foot-and-mouth disease in Great Britain demonstrates the devastating effect microbes can have on livestock and the consequent effect on food supply and economies. Rogue states have actively explored both animal and plant pathogens as weapons.\nLessons learned from the Persian Gulf War highlighted the need for biological warfare agent detectors and the subsequent solutions improved capability on the battlefield. However, other biological hazard (“biohazard”) surveillance deficiencies were soon recognized in the aftermath of conflict. “Gulf War Syndrome” and other ailments suffered by military personnel revealed a need for compact diagnostic tools with integrated sample-processing and detection capabilities to quickly identify disease-causing agents on and off the battlefield. In 1998, a consolidated approach was begun (at the Army Medical Institute for Infectious Diseases) to develop medical diagnostic systems using a common platform for biohazard identification entitled “The Common Diagnostic Systems for Biological Threats and Endemic Infectious Diseases.” Research encompassed development of rapid sample-processing methods, identification technologies, reagents and size reduction of laboratory analysis platforms.\nIn 2002, the Department of Defense (DoD) defined a new approach to a common medical test platform for identifying biological warfare agents and pathogens of operational concern. The Joint Biological Agent Identification and Diagnostic System (JBAIDS) exemplifies this approach. JBAIDS will be configured to support reliable, fast, and specific identification of biological agents from a variety of clinical specimens and environmental sources. JBAIDS will enhance force protection by providing commanders with information to determine actions to protect against and avoid contamination and to restore operations following an attack. JBAIDS information will aid medical personnel in determining appropriate treatment, effective preventive medical measures, and medical prophylaxis in response to the presence of biological agents. Required to combat the threat of biological attack faced by U.S. forces deployed worldwide, JBAIDS will also improve protection against endemic infectious diseases, thereby filling a need identified during the Persian Gulf War for a compact diagnostic identification tool. Today's global military mission, with ongoing operations in war-torn locations teeming with infectious diseases, demands a readily accessible, far-forward biological agent identification capability. This is critical to maintaining troop readiness, quickly determining patient treatment, disposition (for example, quarantine and medical evacuation), and protecting the homeland population from infections acquired by the military, from bioterrorism, and from emerging disease threats.\nThe DoD has addressed the biological threat in the context of the battlefield. However, biological threat reduction in the civilian population context is different. For example, the average civilian is not trained or equipped for response, the public health system is not supported with the kind of central command and control systems associated with the military, different requirements exist on sensitivity and different levels of tolerance for false positives and false negatives, and there is a need for dealing with a broader set of potential agents. Also, much higher sensitivity is required for Counter Terrorism (CT) detection, raising substantial technology challenges and the need to assess background interferences that may be more significant for low-level detection and monitoring schemes.\nThe urgent need for improved biohazard surveillance capability was also recognized and described for the first time in other Government agencies during this period. For example, the United States Postal Service has developed a Biohazard Detection System (BDS) using proven technology to implement early identification of anthrax. The BDS unit consists of an air-collection hood, a cabinet where the collection and analysis devices are housed, a local computer network connection, and a site controller (a networked computer). All BDS processes are automated. The equipment continuously collects air samples from mail canceling equipment while the canceling operation is underway. The air collection hood is installed over the canceling equipment at the very first pinch point in the mail processing operation where it absorbs and concentrates airborne particles into a sterile water base. This creates a liquid sample that is injected into a cartridge. An automated polymerase chain reaction (PCR) test is performed on the liquid sample using sophisticated DNA matching to detect the presence of anthrax (Bacillus anthracis). The test sample is compared to a template for the anthrax DNA sequence for a match. The system concentrates air samples for a one-hour period followed by the PCR test that takes approximately 30 minutes. The BDS is simultaneously concentrating particles for the next sample while the PCR test is performed for the previous sample. So while the first result requires approximately 1½ hours, subsequent results are obtained every hour. Upon detection of a DNA match, the BDS computer network conveys that information to the site controller computer. Local management is notified directly by on-site BDS personnel and also by multiple forms of electronic communication from the BDS site controller. The emergency action plan is activated, the facility's building alarm is sounded and everyone in the building is evacuated. Disadvantageously, the BDS is not adapted for identifying biohazards other than the anthrax spore.\nPractitioners in the art have proposed various solutions to the sampling, detection, analysis, identification and reporting problems associated with the biohazard surveillance requirement. For example, in U.S. Pat. No. 5,895,922, Ho describes a process and apparatus for detection of viable and potentially hazardous biological particles that may be dispersed in an airstream. Ho teaches a method for directing each of the contained particles along a linear path through air, in a sequential manner, and sampling them for determination of their size, whether they are biological and viable, and whether they are present in concentrations greater than background levels. The particle size identifies the particles as respirable or not and the particles are characterized as biological and viable by subjecting each particle in turn, to 340 nm, ultraviolet laser light and looking for the emission of fluorescence, which is typically emitted from bacteria or bacterial spores. Fluorescence detected in the 400–540 nm range signals the presence of nicotinamide adenine dinucleotide hydrogen, which is indicative of biological activity or viability. Ho's apparatus is compact, and power-efficient because he uses a solid state, ultraviolet laser that is actuated only when the particle is passing the laser and only if it is deemed to be a biologically viable candidate, but it is disadvantageous for use in a remote automated surveillance station.\nIn U.S. Pat. No. 6,266,428, Flanigan discloses a system and method for remote detection of hazardous vapors and aerosols by means of two differential spectral signature spectra taken in the field of view at a low spectral resolution. A first linear discriminant optimized for the low spectral resolution is applied to the first spectrum to obtain a first response, and a hazardous cloud is detected automatically in accordance with the first response. A second differential spectral signature spectrum is taken in the field of view at a higher spectral resolution and a second linear discriminant optimized for the higher spectral resolution is applied to the second spectrum to obtain a second response, which is formed into a false-color image and displayed to an operator. The operator discriminates the hazardous cloud in accordance with the image. The first and second linear discriminants can be formed by linear programming. Flanigan's system is disadvantageous for use in a remote automated surveillance station.\nIn U.S. Pat. No. 6,317,080, Baxter discloses a method of tracking airborne substances including the steps of detecting the presence of one or more airborne substances and releasing a tracking balloon into the path of the one or more airborne substances, the tracking balloon having a transmission means and a global positioning means adapted to communicate the latitude and longitude coordinates of the tracking balloon whereby the latitude and longitude coordinates of the tracking balloon are representative of the latitude and longitude of the one or more airborne substances previously detected. Baxter neither considers nor suggests solutions to the remote automated surveillance problem.\nIn U.S. Pat. No. 6,490,530, Wyatt discloses an aerosol hazard classification and early warning network that includes a large number of remote detector and analysis units, which are deployed throughout a region under surveillance for a potentially hazardous aerosol intrusion. Such aerosol threats may originate from fires, volcanic eruptions, or overt releases of biological and chemical agents dispersed in aerosol form. Among the former are the characteristic toxic aerosols released during refinery fires or explosions. The latter biological agents include bacterial spores, lyophilized bacterial cells, and virus preparations, whereas chemical agents might include various forms of nerve gasses and other anti-personnel gasses such as mustard, all commonly deployed in aerosol form. Each detector station contains an aerosol handling unit that samples and transfers ambient aerosol particles one-at-a-time through a light scattering chamber where each such particle is constrained to pass through a fine laser beam producing, thereby, an outgoing scattered light wave. The scattering chamber contains a plurality of scattered light detectors arranged to accept light scattered into different angular locations. The signals detected at each detector position are processed by a corresponding digital signal processing chip with the resulting set of digitized signals being transferred to an on-board central processing unit (CPU). The CPU analyzes the set of light scattering signals and identifies or otherwise characterizes each particle. The classification data are then stored and, on preprogrammed command, telemetered to a remote “central station” by means of an on-board telemetry unit. The central station analyzes the sets of data received from all the detector stations and then instructs, as necessary, selected detector stations via telemetric means to change their sampling and telemetry rates. As soon as sufficient data are available, the central determines the presence, threat, extent, and progress of the aerosol cloud. These factors are then telemetrically transmitted by means of alarms and warnings sent to potentially threatened regions. Although Wyatt's system is well-adapted to remote surveillance and he teaches the use of fluorescence to identify biological compounds, his “light scattering” data are adapted to characterizing and counting particles in an aerosol and Wyatt doesn't consider the rapid and automated identification of biological particles.\nIn U.S. Pat. No. 6,532,067, Chang, et al. describe a method for fluorescence probing of particles flowing in a fluid, including steps of defining a trigger volume in the fluid by intersecting a plurality of substantially orthogonally aimed trigger laser beams, each of a different wavelength, detecting light scattered from the vicinity of the trigger volume by a plurality of particle detectors each sensitive to a wavelength corresponding to the wavelength of a trigger laser beam, probing the particles with a pulsed laser triggered by the particle detectors, collecting fluorescence emitted from the particle in a detection volume and focusing it in a detection region, detecting the fluorescence focused in the detection region. Chang, et al. neither consider nor suggest solutions to the remote automated surveillance problem.\nIn U.S. Pat. No. 6,613,571, Cordery et al. disclose a method and system for detecting biological and chemical hazards in mail using predetermined descriptions of the hazards sought but they neither consider nor suggest solutions to the remote automated surveillance problem. In U.S. Pat. No. 6,656,253, Willey, et al. disclose a dynamic electrostatic filter apparatus for purifying air using electrically charged liquid droplets but they neither consider nor suggest solutions to the remote automated surveillance problem. In U.S. Pat. No. 6,664,550, Rader et al. describe an aerosol lab-on-a-chip (ALOC) that integrates one or more of a variety of particle collection, classification, concentration (enrichment), an characterization processes onto a single substrate or layered stack of such substrates. By mounting a UV laser diode laser light source on the substrate, or substrates tack, so that it is located down-stream of the sample inlet port and at right angle the sample particle stream, the UV light source can illuminate individual particles in the stream to induce a fluorescence response in those particles having a fluorescent signature such as biological particles, some of said particles. An illuminated particle having a fluorescent signal above a threshold signal may trigger a sorter module that separates that particle from the particle stream. But Cordery et al. consider the process control and particle stream separation problems and neither consider nor suggest solutions to the remote automated surveillance problem.\nIn view of the recent terrorism-related security requirements mentioned above, there is a clearly-felt need in the art for a robust (military-hardened) miniaturized remote system for the initial detection, localized analysis and reporting of the presence of biohazards. Such a system requires a large number of permanently-deployed remote surveillance stations each of which can operate independently and without human intervention. Such stations must be adapted for accepting updated detection information from a remote control center to permit adaptation to global changes in the threat environment, for example.\nThese unresolved problems and deficiencies are clearly felt in the art and are solved by this invention in the manner described below."}
{"text": "1. Field of the Invention\nThis invention relates to a display or printing device which is capable, by selective energization of a number of segments, of displaying Arabic letters, numerals and words as well as European numerals.\nThere is a strong need for such a device in the manufacture of such apparatus as small computers, portable testing or communication equipment, etc. where production is suited to markets of the Arab people and where the requirements of relatively small display area and/or low production cost make the use of said device highly desirable.\n2. Description of the Prior Art\nThe main problem in developing such a device is that it is more difficult to devise an arrangement of a limited number of segments which can be selectively operated to display Arabic words than it is for languages using the Latin alphabet. In the case of the latter said languages the starburst method has been successfuly and commonly employed (see, for example, GB No. 2001468) as disclosed, for example, in United Kingdom Pat. No. GB 2,001,468 B. Among the reasons for the aforementioned difficulty is that Arabic alphabetical characters do not lend themselves to simple goemetrical representation as readily as Latin letters. Another reason for the difficulty is that most Arabic letters take on two and sometimes three character forms, depending on their position in the printed word."}
{"text": "The alkylation of benzene with ethylene over a molecular sieve catalyst is a well-known procedure for the production of ethylbenzene. Typically, the alkylation reaction is carried out in a multistage reactor involving the interstage injection of ethylene and benzene to produce an output from the reactor that involves a mixture of monoalkyl and polyalkylbenzene. The principal monoalkylbenzene is, of course, the desired ethylbenzene product. Polyalkylbenzenes include diethylbenzene, triethylbenzene, and xylenes.\nIn many cases, it is desirable to operate the alkylation reactor in conjunction with the operation of a transalkylation reactor in order to produce additional ethylbenzene through the transalkylation reaction of polyethylbenzene with benzene. The alkylation reactor can be connected to the transalkylation reactor in a flow scheme involving one or more intermediate separation stages for the recovery of ethylene, ethylbenzene, and polyethylbenzene.\nTransalkylation may also occur in the initial alkylation reactor. In this respect, the injection of ethylene and benzene between stages in the alkylation reactor not only results in additional ethylbenzene production but also promotes transalkylation within the alkylation reactor in which benzene and diethylbenzene react through a disproportionation reaction to produce ethylbenzene.\nVarious phase conditions may be employed in the alkylation and transalkylation reactors. Typically, the transalkylation reactor will be operated under liquid phase conditions, i.e., conditions in which the benzene and polyethylbenzene are in the liquid phase, and the alkylation reactor is operated under gas phase conditions, i.e., pressure and temperature conditions in which the benzene is in the gas phase. However, liquid phase or critical phase conditions can be used where it is desired to minimize the yield of undesirable by-products from the alkylation reactor.\nIt is a continuing goal of the industry to find and use catalysts that give improved activity and selectivity."}
{"text": "Atom probe is an analytical method capable of directly observing atomic arrangement/composition distribution of a leading end of a sample that has been processed so as to be needle-shaped, at an atomic scale. In the atom probe, a high direct current voltage is applied so as to cause the leading end of the needle-shaped sample to generate a high electric field. A pulse voltage is applied or a pulse laser is irradiated to the leading end so that field evaporation of an atom belonging to a first layer of a surface is induced. Then, mass of an ion that has field-evaporated is time-of-flight-measured so that a type of an element can be determined. Since the field evaporation progresses every atomic layer, atom probe analysis has resolution in a depth direction at an atomic level (refer to PTL 1).\nA peculiar high electric field in an order of 1010 (V/m) due to a metal is required in order to evaporate the metal by an electric field. In the atom probe analysis, a sample should be a needle-shaped sample having a leading end diameter of 100 nm or less in order to achieve the high electric field. In a case where a portion to be observed, such as a bulk sample, is uniformly present over an entire sample, an electric-field polishing method is used for processing of making a leading end of the sample to be acute. In contrast, since a region to be analyzed should be positioned at a leading end of a needle in order to perform the atom probe analysis to a specific region of a semiconductor, an interface, or the like, processing for a needle shape is performed using a focused ion beam (hereinafter, referred to as the “FIB”).\nUsing the FIB enables easy processing of a sample in a needle-shape. However, only the FIB cannot often ascertain whether a leading end of the needle-shaped sample that has been completed actually includes a region to be analyzed. Thus, a sample that has been completed or a leading end of a needle-shaped sample that has been being processed, is sometimes observed by a transmission electron microscope (hereinafter, referred to as a “TEM”). In this case, processing and observation are alternately repeated in order to position a region to be analyzed at a leading end of a needle. Thus, a risk of damaging the sample increases. One holder is preferably shared and repeatedly used between the FIB and the TEM for purposes of a reduction of the risk that the sample is damaged and efficiency of work. PTL 2 describes a sample holder including an atmosphere-isolation mechanism that prevents a sample from altering by an atmospheric effect during a movement between an FIB and a TEM."}
{"text": "The present invention relates in general to techniques for obtaining improved images by compensation methods. In particular it relates to a method for compensating defective pixels in a Spatial Light Modulator (SLM), used in optical lithography. It also relates to an apparatus for patterning a work piece comprising such a method and a method for detecting defective pixels.\nLithographic production is useful for integrated circuits, masks, reticles, flat panel displays, micro-mechanical or micro-optical devices and packaging devices, e.g. lead frames and MCM\"\"s. Lithographic production may involve an optical system to image a master pattern from a computer-controlled reticle onto a workpiece. A suitable workpiece may comprise a layer sensitive to electromagnetic radiation, for example visible or non-visible light. An example of such a system is described in WO 9945435 with the same inventor and applicant as the present invention.\nSaid computer controlled reticle may be a Spatial Light Modulator (SLM) comprising a one or two dimensional array or matrix of reflective movable micro mirrors, a one or two dimensional array or matrix of transmissive LCD crystals, or other similar programmable one or two dimensional arrays based on gratings effects, interference effects or mechanical elements (e.g., shutters).\nIn general, these computer controlled reticles may be used for the formation of images in a variety of ways. These reticles, such as an SLM, include many modulating elements or pixels, in some instances million or more pixels. For example a problem with Spatial Light Modulators is that one or a plurality of pixels in a given SLM may be defective, i.e. they may not respond to a control signal as intended.\nThese defective pixels in a computer controlled reticle may limiting the resolution and accuracy available for their use in optical imaging; e.g., the production of printed patterns on a workpiece may be limited as to its line widths and accuracy.\nTherefore, there is a need in the art for a method, which effectively and precisely finds and compensates for defective pixels in the SLM.\nIn view of the foregoing background, the compensation for defective pixels in the SLM, such as for example a mirror elements stuck in a specific position, is useful to form images having sub micron line widths with tolerances approaching 5 nm.\nAccordingly, it is an object of the present invention to improve the images formed using spatial light modulators by providing an improved method for the compensation of defective pixels.\nIn a first embodiment, the invention provides a method for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation. Said method comprising the actions of providing a source for emitting electromagnetic radiation, illuminating by said radiation said SLM having a plurality of modulating elements (pixels), projecting in a writing pass an image of said modulator on said work piece, and performing a compensation for defective pixels in at least one other writing pass.\nIn another embodiment of the invention said electromagnetic radiation is a pulsed laser source.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a single compensating pixel in another writing pass.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a plurality of compensating pixel in another writing pass.\nIn another embodiment of the invention said SLM is illuminated by a radiation dose in the different writing passes.\nIn another embodiment of the invention said SLM is illuminated by different radiation intensities in the different writing passes.\nIn another embodiment of the invention said SLM is a transmissive Spatial Light Modulator.\nIn another embodiment of the invention said SLM is a reflective Spatial Light Modulator.\nIn another embodiment of the invention the pixels in said SLM is operated in an analog mainer.\nThe invention relates also to a method for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation. Said method comprising the actions of, providing a source for emitting electromagnetic radiation, illuminating by said radiation said SLM having a plurality of modulating elements (pixels), projecting an image of said SLM on the detector arrangement for measuring a dose of radiation, and performing a compensation of said defective pixel by at least one of the most adjacent pixels in the SLM.\nIn another embodiment of the invention said compensation is performed by assigning said at least one of the most adjacent pixels by a value given by subtraction of an intended pixel value by a actual pixel value divided by the number of most adjacent pixels used for compensation.\nThe invention relates also to a method for compensating the impact of at least one defective pixel in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered at least partially with a layer sensitive to electromagnetic radiation. Said method comprising the actions of, setting the pixels in said SLM in a predetermined state, illuminating by a radiation source said SLM, projecting the image of the SLM onto the detector arrangement that measures dose of the SLM pixels, identifying defective pixels, and performing a compensation for said defective pixels in at least one writing pass.\nThe invention relates also to a method for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation. Said method comprising the actions of, providing a source for emitting electromagnetic radiation, illuminating by said radiation said SLM having a plurality of modulating elements (pixels), projecting in a first writing pass an image of said modulator on said work piece using a first set of pixels in said SLM, performing a pre compensation for defective pixels in at least one subsequent writing pass in at least one prior writing pass, and projecting in at least a second writing pass said image of said modulator on said work piece using at least a second set of pixels in said SLM.\nIn another embodiment of the invention, said method further comprising the action of, performing a post compensation for defective pixels in at least one prior writing step in at least one subsequent writing pass.\nIn another embodiment of the invention a post compensation for defective pixels in at least one prior writing step in at least one subsequent writing pass is performed instead of said pre compensation.\nIn another embodiment of the invention said electromagnetic radiation is a pulsed laser source.\nIn another embodiment of the invention, said method further comprising the action of, including at least one pixel in said first set of pixels in said at least second set of pixels.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a single compensating pixel in another writing pass.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a plurality of compensating pixels in another writing pass.\nIn another embodiment of the invention said SLM is illuminated by the same radiation dose in different writing passes.\nIn another embodiment of the invention said SLM is illuminated by different radiation dose in different writing passes.\nIn another embodiment of the invention said SLM is a transmissive Spatial Light Modulator.\nIn another embodiment of the invention said SLM is a reflective Spatial Light Modulator.\nIn another embodiment of the invention the pixels in said SLM is operated in an analog manner.\nIn another embodiment of the invention an image of said pixels in said first writing pass is displaced in said SLM relative said image of said pixels in said second writing pass by one or a plurality of pixels.\nIn another embodiment of the invention an image of said pixels in said first writing pass is displaced on said workpiece relative said image of said pixels in said second writing pass by at least a fraction of a pixel.\nIn another embodiment of the invention said first set of pixels belongs to a first SLM and said second set of pixels belong to a second SLM.\nIn another embodiment of the invention said first and second SLMs are illuminated simultaneously.\nIn another embodiment of the invention said first and second SLMs are illuminated by different radiation intensities.\nThe invention relates also to an apparatus for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation, comprising a source for emitting electromagnetic radiation, a projection system to project in a first writing pass an image of said modulator on said work piece using a first set of pixels in said SLM, means for performing a pre compensation of defective pixels in at least one subsequent writing pass in at least one prior writing pass, a projection system to project in at least a second writing pass said image of said modulator on said work piece using at least a second set of pixels in said SLM, means for performing a post compensation to of defective pixels in at least one prior writing pass in at least one latter writing pass.\nIn another embodiment of the invention said electromagnetic radiation is a pulsed laser source.\nIn another embodiment of the invention at least one pixel in said first set of pixels is included in said at least a second set of pixels.\nIn another embodiment of the invention said projection system to project in at least a second writing pass comprises, said SLM reprogrammed with the image to be written on said work piece with said at least second set of pixels, a movable stage upon which stage said work piece is arranged in order to illuminate the same feature on said work piece with said at least second writing pass as said first writing pass.\nIn another embodiment of the invention said movable stage is moved the length of N SLM pixels.\nIn another embodiment of the invention said stage is moved along a row of pixels.\nIn another embodiment of the invention said movable stage is moved along a column of pixels.\nIn another embodiment of the invention said movable stage is moved along both a row of pixels and a column of pixels.\nIn another embodiment of the invention said movable stage is moved the length of N SLM pixels plus a fraction of a SLM pixel.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a single compensating pixel in another writing pass.\nIn another embodiment of the invention a single defective pixel in one writing pass is compensated with a plurality of compensating pixel in another writing pass.\nIn another embodiment of the invention said SLM is illuminated by a same radiation dose in the different writing passes.\nIn another embodiment of the invention said SLM is illuminated by different radiation intensities in the different writing passes.\nIn another embodiment of the invention said SLM is a transmissive Spatial Light Modulator.\nIn another embodiment of the invention said SLM is a reflective Spatial Light Modulator.\nIn another embodiment of the invention the pixels in said SLM is operated in an analog manner.\nThe invention relates also to an apparatus for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) when creating a pattern of the SLM on a work piece covered with a layer sensitive to electromagnetic radiation, comprising a source for emitting electromagnetic radiation, a projection system for illuminating said SLM, having a plurality of modulating elements (pixels), by said radiation and projecting in a writing pass an image of said modulator on said work piece, the detector arrangement (65) for measuring the dose of pixels from the image of the SLM and a computer (66) for performing a compensation for defective pixels in at least one other writing pass out of said image on said detector (65).\nThe invention relates also to an apparatus for compensating the impact of at least one defective pixel with a known position in a spatial light modulator (SLM) (30) when creating a pattern of the SLM (30) on a work piece (60) covered with a layer sensitive to electromagnetic radiation, comprising a source for emitting electromagnetic radiation, a projection system for illuminating said SLM (30), having a plurality of modulating elements (pixels), by said radiation and projecting in a writing pass an image of said modulator (30) on said work piece (60), the detector arrangement (65) for measuring the dose of pixels from the image of the SLM, and a computer (66) for performing a compensation for defective pixels (110) by using at least one of the most adjacent pixels (111, 112, 113, 114, 115, 116, 117, 118) to said defective pixel (110).\nIn another embodiment of the invention the pixel intensities are detected by said detector (65) whenever a new work piece (60) is to be patterned.\nThe invention also relates to a method for detecting at least one defective pixel in at least one SLM. Said method comprising the actions of addressing all pixels in said at least one SLM with a first steering signal, illuminating said at least one SLM with electromagnetic radiation, recording an image of said at least one SLM, computing a gradient field of the recorded image, computing a divergence of the gradient field, identifying extreme values from the computed divergence which corresponds to defective pixels.\nIn another embodiment said invention further comprising the actions of addressing all pixels said at least one SLM with a second steering signal, illuminating said at least one SLM with electromagnetic radiation, recording an image of said at least one SLM, computing a gradient field of the recorded image, computing a divergence of the gradient field, identifying extreme values from the computed divergence, where defective pixels corresponds to extreme values from said first steering signal and said second steering signal representing same pixels.\nThe invention also relates to a method for detecting at least one defective pixel in at least one SLM. Said method comprising the actions of addressing all pixels in said at least one SLM with a first steering signal, illuminating said at least one SLM with electromagnetic radiation, recording a first image of said at least one SLM, addressing all pixels in said at least one SLM with a second steering signal, illuminating said at least one SLM with electromagnetic radiation, recording a second image of said at least one SLM, computing the difference between said first image and said second image, identifying bad pixels where the computed difference has a local minimum.\nThe invention also relates to a method for detecting at least one defective pixel in at least one SLM Said method comprising the actions of addressing a pattern to said at least one SLM, illuminating said SLM with electromagnetic radiation, recording a first image of said at least one SLM, comparing said recorded image with pattern data at feature edges, identifying bad pixels where the feature edge is moved a predetermined distance.\nIn another embodiment said pattern is a chessboard pattern.\nIn another embodiment said pattern is a pattern with parallel lines.\nIn still another embodiment said method further comprising the action of addressing said pattern with another set of pixels in said at least one SLM, illuminating said SLM with electromagnetic radiation, recording a second image of said at least one SLM, comparing said recorded second image with pattern data at feature edges, identifying bad pixels where the feature edge is moved a predetermined distance.\nIn still another embodiment said method further comprising the action of comparing the feature edge movement in said first image with said second image for identifying bad pixels stuck at intermediate values.\nIn yet another embodiment said different writing passes is performed by means of one SLM.\nIn yet another embodiment said different areas of said SLM are used in the different writing passes.\nIn yet another embodiment said different writing passes is performed by means of a plurality of SLMs.\nIn yet another embodiment said different writing passes is performed by means of one SLM.\nIn yet another embodiment said different areas of said SLM are used in the different writing passes.\nIn yet another embodiment said different writing passes is performed by means of a plurality of SLMs."}
{"text": "Computers and mobile devices, such as cellular phones and personal digital assistants, have become increasingly interconnected due to the widespread availability of wired and wireless connections to communications networks such as the Internet. Even in the earliest days of the ARPANET, users took advantage of such interconnectivity to communicate with one another through early forms of email. As email grew in availability and popularity, email “lists” became a popular tool for broadcasting messages to predefined groups of recipients.\nIn the 1980's, Internet based “newsgroups” emerged in which users could read and respond to discussion threads revolving around a wide variety of predefined categories. Newsgroups are generally readable and updatable by anyone with the equipment to access them, since access to newsgroups is generally not restricted on a per-newsgroup or per-user basis. World Wide Web based discussion groups (i.e., also known as clubs) have also provided a way for groups of people to associate around a topic. Through the use of web server programming, the idea of discussion groups and discussion threads has been extended to provide users with the ability to subscribe to secured discussion forums that are, in some cases, moderated by other users.\nAnother variant of Internet based communication forums are the web-based “social network” applications, in which a number of users are able to find each others' accounts and voluntarily become “friends” or “followers” of each other's posted messages. Users generally post brief messages about their status, mood, activities, and such, and their friends and followers can read and optionally reply to those messages. As such, friends may stay abreast of each other's' activities as a tool for maintaining their social bonds."}
{"text": "The invention relates to bird feeders, and particularly feeders designed for relatively small birds, such as finches, chickadees, etc.\nThe instant invention relates to bird feeders of the general type having an elongated, transparent, vertically disposed cylindrical housing preferably constructed of a suitable plastic material, and having adjustable feed ports at the lower extremity thereof, with a peripheral skirt or tray member located beneath the feed ports for collecting seed that spills outwardly through the ports and at the same time for supporting the type of small birds for which this feeder is designed. Feeders of this general type are known in the prior art, but assembly and disassembly of such feeders during cleaning and/or repair frequently is a problem. It is therefore an object of my invention to provide such a feeder wherein the component parts can easily and quickly be assembled and disassembled, such assembly and disassembly requiring no special tools other than an ordinary screwdriver or the like. In fact, assembly and disassembly of the bird feeder of the present invention, not including the top cap or cover thereof, is achieved by manipulating a single screw member, which is easily and conveniently accessible."}
{"text": "The present invention relates to interconnect structures for high-speed microprocessors, application specific integrated circuits (ASICs), and other high-speed integrated circuits (ICs). More particularly, the present invention provides an interconnect structure of the dual damascene-type which includes a low dielectric constant (i.e., low-k) dielectric having Cu regions formed therein. The Cu interconnect structures of the present invention have enhanced circuit speed and reduced fabrication cost. Moreover, the inventive interconnect structures contain an inventive diffusion barrier layer that is used as a post-CMP (chemical-mechanical polishing) cap.\nMany low-k dielectrics having a dielectric constant, k, of about 3.5 or less, plus Cu interconnect structures are known in the art; See, for example, R. D. Goldblatt, et al., xe2x80x9cA High Performance 0.13 xcexcm Copper BEOL Technology with Low-K Dielectricxe2x80x9d, Proceedings of the International Interconnect Technology Conference, IEEE Electron Devices Society, Jun. 5-7, 2000 pgs. 261-263. Such prior art interconnect structures include inorganic, organic or mixed dielectric materials as the interlevel or intralevel dielectric. It is widely accepted that dual-damascene structures are lower cost than single damascene or subtractive metal structures.\nTypically, there are two essential problems associated with prior art Cu interconnect structures which include: Cu migration out of the conductive regions into the surrounding dielectric which causes a reliability problem with the structure; and increasing cost of fabricating the interconnect structure due to increasing number of processing steps and expensive tools. A discussion concerning the problem of Cu migration in Cu interconnect structures is found, for example, in A. Loke, et al., xe2x80x9cEvaluation of Cu Penetration in Low-k Polymer Dielectrics by Bias-Temperature Stressxe2x80x9d, Mat. Res. Soc. Symp. Proc., Vol. 565 (1999) p. 173.\nCommon solutions to the above problems add additional processing steps, thus further increasing the cost to fabricate the desired low-k plus Cu interconnect structure. Furthermore, the post-CMP cap layers in prior art dual damascene structures are made using vacuum-based plasma-enhanced chemical vapor deposition (PECVD) tools that are costly to purchase and maintain.\nCo-assigned U.S. application Ser. No. 09/371,340 filed Aug. 10, 1999 provides an interlayer dielectric that is capable of reducing or eliminating Cu ion migration. Specifically, the interlayer dielectric of the \"\"340 disclosure comprises a dielectric material having a dielectric constant of 3 or less and an additive. The additive employed therein is capable of binding Cu ions, is soluble in the dielectric material and is substantially, uniformly distributed throughout the dielectric. This disclosed interlayer dielectric is not used as a spin-on dielectric in the \"\"340 disclosure.\nIn view of the above-mentioned problems in prior art processes of fabricating dual damascene-type structures, there is a continued need for providing a new and improved method for forming such interconnect structures which substantially minimizes Cu migration into the surrounding dielectric layers, without adding extra processing steps and cost in fabricating the interconnect structure.\nOne object of the present invention is to provide an interconnect structure which includes a new diffusion barrier layer (i.e., post-CMP cap) that substantially minimizes the migration of Cu into the surrounding low-k dielectrics.\nA further object of the present invention is to provide a Cu interconnect structure of the dual damascene-type that is made without any vacuum-based PECVD tools.\nA still further object of the present invention is to provide a Cu interconnect structure that has an effective low dielectric constant (on the order of about 3.5 or less), which is easy to manufacture and does not include extra processing steps that would increase the cost of fabricating the same.\nThese and other objects and advantages are achieved in the present invention by providing a diffusion barrier layer (i.e., post-CMP cap) that is based on a spin coated dielectric layer. Specifically, the diffusion barrier layer of the present invention (hereinafter referred to as spin-on cap) comprises a low-k dielectric (on the order of about 3.5 or less) that includes at least one additive. The at least one additive employed in the present invention is capable of binding Cu ions, and is soluble in the spun-on low-k dielectric. The spin-on cap of the present invention may further include a spun-on low-k (on the order of 3.5 or less) reactive-ion etch (RIE) stop layer. Spin-on caps containing a bilayer of low-dielectric plus at least one additive and low-k RIE stop layer are preferred. It is noted that the inventive spin-on cap does not significantly increase the effective dielectric constant of the interconnect structure and does not add additional cost to the fabrication of the interconnect structure since a single deposition tool, i.e., spin coating tool, is employed.\nOne aspect of the present invention relates to a Cu interconnect structure which includes the inventive spin-on cap as a post-CMP cap. Specifically, the inventive Cu interconnect structure comprises:\na substrate having a patterned low-k interlevel dielectric formed thereon, said patterned low-k interlevel dielectric having an effective dielectric constant of about 3.5 or less;\nCu conductive regions formed within said patterned low-k interlevel dielectric;\na polish stop layer on surfaces of said patterned low-k interlevel dielectric not containing said Cu conductive regions; and\na spin-on cap present on said Cu conductive regions and said polish stop layer, wherein said spin-on cap comprises a low-k dielectric having a dielectric constant of about 3.5 or less and at least one additive, wherein said at least one additive is capable of binding Cu ions, and is soluble in the spun-on low-k dielectric.\nIn one embodiment of the present invention, the polish stop layer is not present on the patterned low-k interlevel dielectric. Instead, in this embodiment, the spin-on cap is formed on the patterned low-k interlevel dielectric.\nAnother aspect of the present invention relates to a method of fabricating the above mentioned Cu interconnect structure. Specifically, the method of the present invention comprises the steps of:\n(a) forming a low-k interlevel dielectric on a surface of a substrate, said low-k interlevel dielectric having an effective dielectric constant of about 3.5 or less;\n(b) forming a hard mask on said low-k interlevel dielectric, said hard mask including at least a polish stop layer;\n(c) forming an opening in said hard mask so as to expose a portion of said low-k interlevel dielectric;\n(d) forming a trench in said low-k interlevel dielectric using said hard mask as an etch mask;\n(e) filling said trench with at least Cu;\n(f) planarizing said Cu stopping on said polish stop layer; and\n(g) forming a spin-on cap on said Cu and said polish stop layer, wherein said spin-on cap comprises a low-k dielectric having a dielectric constant of about 3.5 or less and at least one additive, wherein said at least one additive is capable of binding Cu ions, and is soluble in the spun-on low-k dielectric.\nThe term xe2x80x9ctrenchxe2x80x9d as used herein denotes a line, via or a combination of a line and via. In one embodiment of the present invention, step (b) above is omitted."}
{"text": "1. Field of the Invention\nThis invention relates to novel acetal terpolymers. More particularly, this invention relates to novel ultraviolet (UV) curable terpolymers of trioxane, 1,3-dioxolane and certain formals of monoethylenically unsaturated aliphatic diols. These terpolymers are non-crystalline at room temperature (about 25.degree. C.) or above by virtue of having higher dioxolane contents than copolymers or terpolymers hitherto contemplated by the prior art, although they can be made to crystallize at temperatures below room temperature. They can be admixed with multifunctional acrylates or similarly - performing multifunctional crosslinking monomers and a photosensitizer or UV initiator and cured to an insoluble, non-tacky, rubbery state using UV radiation. The resulting cured polymeric materials form useful crosslinked films, and can also be cryogenically ground to a suitable small particle size and then blended with conventionally prepared crystalline oxymethylene homo-, co- and terpolymers to improve the impact properties of molded articles made therefrom.\n2. Description of the Prior Art\nCommonly-assigned copending U.S. patent application Ser. No. 07/096,187, filed of even date herewith in the names of George L. Collins, Paul Zema and William Pleban and entitled \"Low Tg Non-Crystalline Acetal Copolymers\", discloses and claims low Tg (glass transition temperature) trioxane/ 1,3-dioxolane copolymers which are also non-crystalline at room temperature or above (although they too can be made to crystallize at temperatures below room temperature) and which have a dioxolane content greater than 65 mol percent and less than about 75 mol percent and an IV (intrinsic viscosity) of from about 1.00 to about 2.3, as measured by standard viscometric measurements, e.g., in o-chlorophenol. These copolymers contain no unsaturated sites for subsequent crosslinking, and are not disclosed as being curable by reaction with multifunctional crosslinking monomers under UV curing conditions.\nTrioxane/1,3-dioxolane/unsaturated diol formal terpolymers are disclosed in U.S. Pat. No. 3,297,647, issued January 10, 1967 to Schott et al. The unsaturated diol formals disclosed in the Schott et al patent include 2-butene-1,4-diol formal, 2-hexane-1,4-diol formal and 2-(2-ethyl)butene-1,4-diol formal (column 2, lines 11-14), the amounts of monomers disclosed are as follows:\n\"(t)he cyclic ether or saturated cyclic formal is advantageously used in an amount of 0.1 to 59.9% by weight, calculated on the total monomer mixture. The formal of the unsaturated diol is advantageously used in an amount of 59.9 to 0.2% by weight, calculated on the total monomer mixture. The trioxane is advantageously used in an amount of 40 to 99.8% by weight, calculated on the total monomer mixture\", PA0 \"The physical properties of the terpolymers can be varied within wide limits and depend, on the one hand, on the nature and concentration of the saturated cyclic formal or saturated cyclic ether and, on the other hand, on the concentration of the formal of an unsaturated cyclic diol. PA0 For example, when 0.1 to 10% by weight, calculated on the total monomer mixture, of unsaturated cyclic formal and 0.1 to 10% by weight, calculated on the total monomer mixture, of saturated cyclic formal or ether are used highly crystalline products are obtained, whereas with 40 to 59.9% by weight, calculated on the total monomer mixture of unsaturated cyclic formal or ether, amorphous elastic glass-clear products are obtained. The more voluminous the second comonomer, the lower is the crystallinity and the higher is the elasticity. The decrease in crystallinity can be well measured by means of X-rays. PA0 Low molecular weight terpolymers which constitute waxes or oils can easily be obtained with the use of high concentrations of catalyst, that is about 0.1 to 1% by weight, calculated on the total monomer mixture. PA0 The above statements are intended to indicate the wide limits within which the properties of the terpolymers obtained by the process of the invention may be varied, the incorporation of different comonomers having, of course, different effects on the properties of the terpolymers and the transitions being fluid\", PA0 \"(v)ery interesting properties are imparted to the terpolymers of the invention by the double bonds contained in the main chain, which double bonds enable the terpolymers to be crosslinked by known methods. For example, the terpolymers may be vulcanized by kneading with sulfur\", PA0 \". . . such that at least 30% of the polymer composition remains undissolved when a film thereof of 5-8 mils thickness is immersed in one hundred times its weight of boiling dimethylformamide for two minutes\", PA0 \". . . at least one covalently bonded linkage [is present] between catenarian carbon atoms of different polyoxymethylene chains for each for polyoxymethylene polymer molecules\",\ncolumn 2, lines 17-24, and the patentees had this to say about the interrelationships between monomer contents and polymer properties:\ncolumn 2, lines 25-52.\nThe Schott et al patent does not disclose terpolymers having 1,3-dioxolane contents in excess of 59.9% by weight, based on the total weight of monomers present, or blends of its terpolymers with crystalline oxymethylene polymers. And while Schott et al do disclose that their olefinic unsaturation-containing terpolymers can be crosslinked:\ncolumn 2, lines 53-58, there is neither a disclosure of crosslinking with multifunctional crosslinking monomers nor a disclosure of UV curing.\nU.S. Pat. No. 3,215,671, issued Nov. 2, 1965 to Melby, discloses crosslinking oxymethylene homopolymers using, in one embodiment, from 0.5 to 20% by weight of a polyunsaturated compound (multifunctional acrylates are disclosed) and a photoinitiator under UV light. See, e.g., column 2, lines 21-25 and 63-72. The crosslinked product is characterized as having an extent of crosslinking:\ncolumn 2, lines 30-39, or such that:\ncolumn 2, lines 40-44. Melby's crosslinked oxymethylene homopolymers appear to be crystalline materials; see from column 5, line 67 to column 8, line 10. In no case are they characterized as noncrystalline."}
{"text": "The utilization of an electric current to produce an ameliorative effect on a bodily tissue has a long history, reportedly extending back to the ancient Greeks. The effects on bodily tissue from an applied electric current, and thus the dividing line between harmful and curative effects, will vary depending upon the voltage levels, current levels, the length of time the current is applied, and the tissue involved. One such effect resulting from the passage of an electric current through tissue is heat generation.\nBody tissue, like all non-superconducting materials, conducts current with some degree of resistance. This resistance creates localized heating of the tissue through which the current is being conducted. The amount of heat generated will vary with the power P deposited in the tissue, which is a function of the product of the square of the current I and the resistance R of the tissue to the passage of the current through it (P=I2R.).\nAs current is applied to tissue, then, heat is generated due to the inherent resistance of the tissue. Deleterious effects in the cells making up the tissue begin to occur at about 42° Celsius. As the temperature of the tissue increases due to heat generated by the tissue's resistance, the tissue will undergo profound changes and eventually, as the temperature becomes high enough, that is, generally greater than 45° C., the cells will die. The zone of cell death is known as a lesion and the procedure followed to create the lesion is commonly called an ablation. As the temperature increases beyond cell death temperature, complete disintegration of the cell walls and cells caused by boiling off of the tissue's water can occur. Cell death temperatures can vary somewhat with the type of tissue to which the power is being applied, but generally will begin to occur within the range of 45° to 60° C., though actual cell death of certain tissue cells may occur at a higher temperature.\nIn recent times, electric current has found advantageous use in surgery, with the development of a variety of surgical instruments for cutting tissue or for coagulating blood. Still more recently, the use of alternating electric current to ablate, that is, kill, various tissues has been explored. Typically, current having a frequency from about 3 kilohertz to about 300 gigahertz, which is generally known as radiofrequency or radiofrequency (RF) current, is used for this procedure. Destruction, that is, killing, of tissue using an RF current is commonly known as radiofrequency ablation. Often radiofrequency ablation is performed as a minimally invasive procedure and is thus known as radiofrequency catheter ablation because the procedure is performed through and with the use of a catheter. By way of example, radiofrequency catheter ablation has been used to ablate cardiac tissue responsible for irregular heart beats or arrythmias.\nThe prior art applications of current to tissue have typically involved applying the current using a “dry” electrode. That is, a metal electrode is applied to the tissue desired to be affected and a generated electric current is passed through the electrode to the tissue. A commonly known example of an instrument having such an operating characteristic is an electrosurgical instrument known as a “bovie” knife. This instrument includes a cutting/coagulating blade electrically attached to a current generator. The blade is applied to the tissue of a patient and the current passes through the blade into the tissue and through the patient's body to a metal base electrode or ground plate usually placed underneath and in electrical contact with the patient. The base electrode is in turn electrically connected to the current generator so as to provide a complete circuit.\nAs the current from the bovie knife passes from the blade into the tissue, the resistance provided by the tissue creates heat. In the cutting mode, a sufficient application of power through the bovie knife to the tissue causes the fluid within the cell to turn to steam, creating a sufficient overpressure so as to burst the cell walls. The cells then dry up, desiccate, and carbonize, resulting in localized shrinking and an opening in the tissue. Alternatively, the bovie knife can be applied to bleeding vessels to heat and coagulate the blood flowing therefrom and thus stop the bleeding.\nAs previously noted, another use for electrical instruments in the treatment of the body is in the ablation of tissue. To expand further on the brief description given earlier of the ablation of cardiac tissue, it has long been known that a certain kind of heart tissue known as sino-atrial and atrio-ventricular nodes spontaneously generate an electrical signal that is propagated throughout the heart along conductive pathways to cause it to beat. Occasionally, certain heart tissue will “misfire,” causing the heart to beat irregularly. If the errant electrical pathways can be determined, the tissue pathways can be ablated and the irregular heartbeat remedied. In such a procedure, an electrode is placed via a catheter into contact with the tissue and then current is applied to the tissue via the electrode from a generator of RF current. The applied current will cause the tissue in contact with the electrode to heat. Power will continue to be applied until the tissue reaches a temperature where the heart tissue dies, thereby destroying the errant electrical pathway and the cause of the irregular heartbeat.\nAnother procedure using RF ablation is transurethral needle ablation, or TUNA, which is used to create a lesion in the prostate gland for the treatment of benign prostatic hypertrophy (BPH) or the enlargement of the prostate gland. In a TUNA procedure, a needle having an exposed conductive tip is inserted into the prostate gland and current is applied to the prostate gland via the needle. As noted previously, the tissue of the prostate gland heats locally surrounding the needle tip as the current passes from the needle to the base electrode. A lesion is created as the tissue heats and the destroyed cells may be reabsorbed by the body, infiltrated with scar tissue, or just become non-functional.\nWhile there are advantages and uses for such “dry” electrode instruments, there are also several notable disadvantages. One of these disadvantages is that during a procedure, coagulum—dried blood cells and tissue cells—will form on the electrode engaging the tissue. Coagulum acts as an insulator and effectively functions to prevent current transfer from the blade to the tissue. This coagulum “insulation” can be overcome with more voltage so as to keep the current flowing, but only at the risk of arcing and injuring the patient. Thus, during surgery when the tissue is cut with an electrosurgical scalpel, a build-up of coagulated blood and desiccated tissue will occur on the blade, requiring the blade to be cleaned before further use. Typically, cleaning an electrode/scalpel used in this manner will involve simply scraping the dried tissue from the electrode/scalpel by rubbing the scalpel across an abrasive pad to remove the coagulum. This is a tedious procedure for the surgeon and the operating staff since it requires the “real” work of the surgery to be discontinued while the cleaning operation occurs. This procedure can be avoided with the use of specially coated blades that resist the build up of coagulum. Such specialty blades are costly, however.\nA second disadvantage of the dry electrode approach is that the electrical heating of the tissue creates smoke that is now known to include cancer-causing agents. Thus, preferred uses of such equipment will include appropriate ventilation systems, which can themselves become quite elaborate and quite expensive.\nA further, and perhaps the most significant, disadvantage of dry electrode electrosurgical tools is revealed during cardiac ablation procedures. During such a procedure, an electrode that is otherwise insulated but having an exposed, current carrying tip is inserted into the heart chamber and brought into contact with the inner or endocardial side of the heart wall where the ablation is to occur. The current is initiated and passes from the current generator to the needle tip electrode and from there into the tissue so that a lesion is created. Typically, however, the lesion created by a single insertion is insufficient to cure the irregular heartbeat because the lesion created is of an insufficient size to destroy the errant electrical pathway. Thus, multiple needle insertions and multiple current applications are almost always required to ablate the errant cardiac pathway, prolonging the surgery and thus increasing the potential risk to the patient.\nThis foregoing problem is also present in TUNA procedures, which similarly require multiple insertions of the needle electrode into the prostate gland. Failing to do so will result in the failure to create a lesion of sufficient size otherwise required for beneficial results. As with radiofrequency catheter ablation of cardiac tissue, then, the ability to create a lesion of the necessary size to alleviate BPH symptoms is limited and thus requires multiple insertions of the electrode into the prostate.\nA typical lesion created with a dry electrode using RF current and a single insertion will normally not exceed one centimeter in diameter. This small size—often too small to be of much or any therapeutic benefit—stems from the fact that the tissue surrounding the needle electrode tends to desiccate as the temperature of the tissue increases, leading to the creation of a high resistance to the further passage of current from the needle electrode into the tissue, all as previously noted with regard to the formation of coagulum on an electrosurgical scalpel. This high resistance—more properly termed impedance since typically an alternating current is being used—between the needle electrode and the base electrode is commonly measured by the RF current generator. When the measured impedance reaches a pre-determined level, some prior art generators will discontinue current generation. Discontinuance of the ablation procedure under these circumstances is necessary to avoid injury to the patient.\nThus, a typical procedure with a dry electrode may involve placing the needle electrode at a first desired location; energizing the electrode to ablate the tissue; continue applying current until the generator measures a high impedance and shuts down; moving the needle to a new location closely adjacent to the first location; and applying current again to the tissue through the needle electrode. This cycle of electrode placement, electrode energization, generator shut down, electrode re-emplacement, and electrode re-energization, will be continued until a lesion of the desired size has been created. As noted, this increases the length of the procedure for the patient. Additionally, multiple insertions increases the risk of at least one of the placements being in the wrong location and, consequently, the risk that healthy tissue may be undesirably affected while diseased tissue may be left untreated. The traditional RF ablation procedure of using a dry ablation therefore includes several patient risk factors that both patient and physician would prefer to reduce or eliminate.\nThe therapeutic advantages of RF current could be increased if a larger lesion could be created safely with a single positioning of the current-supplying electrode. A single positioning would allow the procedure to be carried out more expeditiously and more efficiently, reducing the time involved in the procedure. Larger lesions can be created in at least two ways. First, simply continuing to apply current to the patient with sufficiently increasing voltage to overcome the impedance rises will create a larger lesion, though almost always with undesirable results to the patient. Second, a larger lesion can be created if the current density, that is, the applied electrical energy, could be spread more efficiently throughout a larger volume of tissue. Spreading the current density over a larger tissue volume would correspondingly cause a larger volume of tissue to heat in the first instance. That is, by spreading the applied power throughout a larger tissue volume, the tissue would heat more uniformly over a larger volume, which would help to reduce the likelihood of generator shutdown due to high impedance conditions. The applied power, then, will cause the larger volume of tissue to be ablated safely, efficiently, and quickly.\nResearch conducted under the auspices of the assignee of the present invention has focused on spreading the current density throughout a larger tissue volume through the creation, maintenance, and control of a “virtual electrode” within or adjacent to the tissue to be ablated. A virtual electrode can be created by the introduction of a conductive fluid, such as isotonic or hypertonic saline, into or onto the tissue to be ablated. The conductive fluid will facilitate the spread of the current density substantially equally throughout the extent of the flow of the conductive fluid, thus creating an electrode—a virtual electrode—substantially equal in extent to the size of the delivered conductive fluid. RF current can then be passed through the virtual electrode into the tissue.\nA virtual electrode can be substantially larger in volume than the needle tip electrode typically used in RF interstitial ablation procedures and thus can create a larger lesion than can a dry, needle tip electrode. That is, the virtual electrode spreads or conducts the RF current density outward from the RF current source—such as a current carrying needle, forceps or other current delivery device—into or onto a larger volume of tissue than is possible with instruments that rely on the use of a dry electrode. Stated otherwise, the creation of the virtual electrode enables the current to flow with reduced resistance or impedance throughout a larger volume of tissue, thus spreading the resistive heating created by the current flow through a larger volume of tissue and thereby creating a larger lesion than could otherwise be created with a dry electrode.\nWhile the efficacy of RF current ablation techniques using a virtual electrode has been demonstrated in several studies, the currently available instruments useful in such procedures lags behind the research into and development of hoped-for useful treatment modalities for the ablation of soft tissue and malignancies.\nIt would be desirable to have an apparatus capable of creating a virtual electrode for the controlled application of tissue ablating RF electric current to a tissue of interest so as to produce a lesion of desired size and configuration."}
{"text": "Hundreds of man-made satellites are currently in orbit around the Earth. Over the next decade, governments and companies around the globe plan to launch hundreds of new spacecraft for a variety of communications, defense and remote sensing projects. The placement of satellites into Earth orbit can cost many millions of dollars. A conventional launch involves a large multi-stage, single-use rocket to lift a satellite into a geosynchronous orbit.\nA general description of conventional nuclear-propulsion systems may be found in a text entitled A Critical Review of Space Nuclear Power and Propulsion, edited by Mohamed S. El-Genk, which was published by the American Institute of Physics in 1994.\nThe U.S. Departments of Energy and Defense and NASA developed plans for a Generic Flight System for space-based defense systems and NASA exploration missions called SP-100 in the mid-1980's. The SP-100 was designed to supply nuclear-power for military and civilian space systems. This early system was designed as a single-use power stage for a single, permanently attached payload; and was never configured for any on-orbit rendezvous, docking or servicing missions. The SP-100 is described in the SP-100 Technical Summary Report, which was prepared for the U.S. Department of Energy by the Jet Propulsion Laboratory and the California Institute of Technology in September, 1994.\nVarious nuclear electric propulsion systems are described in a publication entitled Nuclear Electric Propulsion, A Summary of Concepts Submitted to the NASA/DoE/DoD Nuclear Electric Propulsion Workshop, which was held in Pasadena, Calif. on 19–22 Jun. 1990.\nThe Aerospace Division of the Olin Corporation proposed a small engine for the small satellite community called the Small Upper Stage (SUS). The SUS was designed to accomplish low Earth orbit transfers, orbit circularizations and plane changes using hydrazine propulsion.\nTRW has patented several methods and apparatus intended for the space transportation market. In U.S. Pat. No. 4,471,926, Steel describes a Transfer Vehicle for Use in Conjunction with a Reusable Space Shuttle. This spacecraft has a propulsion system that uses a low-thrust bi-propellant liquid rocket engine to provide a soft, low-acceleration ascent. In U.S. Pat. No. 4,575,029, Harwood and Love disclose a spacecraft for transporting a payload from a space shuttle in a low altitude parking orbit to an operational orbit. In U.S. Pat. No. 4,943,014, Harwood and Love reveal their “soft ride” method for changing the altitude or position of a spacecraft in orbit using a liquid bi-propellant engine.\nIn U.S. Pat. No. 4,664,344, Harwell describes an apparatus and method of capturing an orbiting spacecraft. This device comprises a relatively small mechanical probe and fixture operated by an astronaut during a spacewalk.\nIn an article entitled Topaz Two Proves to Be a Gem for International Tech Transfer, contained in Technical Applications Report from Ballistic Missile Defense Organization, 1995, thermoionic reactors for space-based power generation are disclosed.\nProspects for Nuclear Electric Propulsion Using Closed-Cycle Magnetohydrodynamic Energy Conversion, by R. Litchford et al. was presented at the 12th Annual Advanced Space Propulsion Workshop in Huntsville, Ala. on 3–5 Apr. 2001.\nJ. Collins et al. disclose a Small Orbit Transfer Vehicle for On-Orbit Servicing and Resupply which was presented at the 15th Annual Utah State University Conference on Small Satellites at Logan, Utah, 13–16 Aug. 2001.\nIn U.S. Pat. No. 4,754,601, Minovitch discloses “a propulsion system for reusable space-based vehicles is presented wherein the propulsive working fluid is atmospheric gas.”\nIn U.S. Pat. No. 5,260,639, De Young et al. describe “a method of supplying power to a device such as a lunar rover located on a planetary surface.”\nIn U.S. Pat. No. 6,213,700, Koppel discloses a “method [which] serves to place a space vehicle, such as a satellite, on a target orbit such as the orbit adapted to normal operation of the space vehicle and starting from an elliptical initial orbit that is significantly different from, and in particular more eccentric than the target orbit.”\nIn U.S. Pat. No. 6,357,700, Provitola describes “an spacecraft/airship, which uses buoyancy and thrusters to ascend into space with lifting gas as propellant or fuel for thrusters, which may be conventional thrusters or electric turbojets or ion thrusters.”\nIn U.S. Pat. No. 5,260,639, Basuthakur et al reveal “a satellite assembly [that] is formed from any number of bus modules which have a substantially common shape and interior space volume.”\nIn U.S. Pat. No. 6,478,257, Oh et al. describe “systems and methods that employ a phase change material to provide thermal control of electric propulsion devices.”\nIn U.S. Pat. No. 3,825,211, Minovitch presents a “space vehicle [which] carries a vaporizable propellant . . . [E]nergy is transmitted to the vehicle while in space by a laser beam originating on the ground or some other body or satellite.”\nThe development of an in-orbit space transportation and rescue vehicle would dramatically reduce the cost of changing the orbital position of a satellite. Such a system would revolutionize the military and commercial space industries, and fill a long-felt need in the telecommunications, direct-broadcast and remote-sensing industries."}
{"text": "Medicament delivery devices are routinely used by persons without formal medical training, i.e. patients where self-management of their condition is increasingly common. These circumstances set a number of requirements for medicament delivery devices of this kind The injector should be robust in construction, yet easy to use in terms of its operation by a user and the manipulation of the parts. In the case of those with diabetes, many users will be of impaired vision and may also be physically infirm. Devices that are too large or cumbersome may therefore prove difficult to use, particularly someone with reduced dexterity.\nU.S. Pat. No. 5,989,221 describes an electronically controlled injection device in which the readying of the device for administering and the subsequent drug delivery therefrom is controlled by an electronic control unit. Specifically, the control unit comprises a position or attitude sensor for transmitting a signal so that the readying of the device cannot take place unless the longitudinal axis of the injection cartridge is in a predetermined direction. This is in the context of removing air or mixing drug components. The control unit is also operative for driving a piston rod a predetermined distance for delivery of a drug dose.\nIn such prior art devices, the user is still required to prime the device after installation of the medicament cartridge. Moreover, such devices may under dose particularly in cases where the medicament cartridge is not properly seated within the device after insertion by the user."}
{"text": "Boats, and particularly pontoon boats, typically include seats. Some seats may be used for lounging and reclining. The present invention relates to reclining seats for use on boats. Embodiments of the present invention may be used on boats including pontoon boats."}
{"text": "Full disk encryption (FDE) technologies are designed to protect the data in case the platform is stolen. Such technologies can be either software-based or hardware-based. These technologies rely on the end-user providing a password on boots from certain states in order to unlock the access to data stored on device. However, FDE protects a computer's data-at-rest only when it is not decrypted yet, e.g. when it is being booted.\nAnother theft protections system is a software-based alerting mechanism. Software-based alerting mechanisms provide an immediate alert capability in order to protect against theft. The problem is that these mechanisms are susceptible to software-based attacks by thieves (e.g., turning off the WIFI radio), simple hardware-based attacks by thieves (e.g., pressing the platform's power button for 4 seconds).\nAnother theft protection system relies on discrete hardware components containing trigger-based alerting mechanisms. An example for this is a disk-on-key like component that gets plugged into the PC. However, this requires an additional plug-in device, and only works when the computer system is already active. In addition, a thief can easily destroy such components while keeping the platform intact, e.g. drown it in a glass of water, or bring down a hammer on it."}
{"text": "It has become well-established in the packaging field to weigh bulk articles in a combinational computerized weigher before packaging, thus to obtain a high degree of accuracy when packaging bulk materials. When the products weighed for packaging comprise rod-like elements such as licorice, a group of pencils, pretzel sticks, bread sticks, and the like, there is a particular problem of how best to handle such elongate products for packaging with a minimum of breakage. The ideal condition is to package these articles as an aligned group so as to consume a minimum amount of packaging material and to make a most compact package for placing subsequently in the larger shipping container or box. In the field it is generally understood that conventional packaging techniques which result in random placement of elongate articles also may require larger volume of packages than if the same number of articles were placed in alignment and then packaged.\nTo facilitate rapid weighing and packaging of elongate articles, machines have been developed which align the articles prior to packaging. U.S. Pat. No. 4,965,984 discloses a method and apparatus for handling elongate elements on an inclined flight-equipped conveyer which discharges into a collector for storage and transfer to a packaging machine. The articles free-fall from the conveyor to the collector.\nU.S. Pat. No. 4,586,313 concerns the packaging of elongated pieces characterized by horizontal, planer series of pans in which French fries and the like are moved from pan to pan and aligned with the longitudinal axes parallel to one other by the action of a plurality of vibrators. When the individual pieces are aligned, they fall into a weighing bucket which has an electronically controlled retention flap that will provide the aligned French fries into a bagging machine when a full charge is received.\nU.S. Pat. No. 4,514,959 discloses a system for aligning and packaging elongated products in two steps. First, the randomly oriented product, such as French fries, are horizontally disposed on a vibrating conveyor. At a distal end of the conveyor, a number of channels are provided which receive the advancing French fries. The channels are configured so that the articles must advance with their elongate dimensions parallel to the channel length and the channels are positioned at a preselected angle with respect to a discharge chute. When the French fry is presented to the end of the channel towards the discharge chute, it is turned further to be parallel with the length of the discharge chute. The French fry then drops through the chute into a weighing device and thence to the packaging machine.\nWe believe it desirable to minimize the distance and time a product is in free-fall in order to materially reduce the amount of product fracture and crumbling and to gain efficiency in handling elongate products.\nWe believe it would be advantageous to provide an apparatus and method for receiving weighed charges of elongated articles and for presenting those articles in as gentle a condition as feasible to the packaging machine and the present invention is focused at that result."}
{"text": "In conventional systems, users have access to a vast amount of media content. A user may search for content that he or she finds interesting by navigating through numerous pages of content. Oftentimes, a user may find the process of searching to be boring or cumbersome, and may become disinterested in consuming media completely. Some conventional systems solve this problem by enabling a user to play games instead of consuming content. However, those games are usually unrelated to media content that the user finds interesting."}
{"text": "With reference to FIG. 1, an exploded view of Taiwan Patent No. 0350530 is disclosed. The wind outlet plate disclosed includes a wind outlet panel A, supports B, a bearing box C and a negative ion generator D.\nFour supports B are installed at the back side of the wind outlet panel A. The supports B serve to support the bearing box C of the negative ion generator D on the wind outlet panel A.\nIt will be appreciated from the prior art that although the primary structure is simple, it is very difficult to detach the structure. Since the structure in the prior art has no filter. Many dusts possibly flow into a room so that the quality of air can not be assured. Moreover, since the prior art has a fixed structure, if a filter is desired to be attached thereon, the wind outlet panel is necessary to be detached. Moreover, since the negative ion generator in the prior art is controlled by a central controlled air conditioner, in order to avoid danger accidents, the air conditioner must be closed for being detached and cleaned.\nAlthough the prior art can provide a preferred air, the maintenance and cleaning works are difficult. For example, closing the central controlled air conditioner will affect the works of other peoples. If the central controlled air conditioner is not closed, it is possible that the electric power may be drained out. Moreover, the effects in dust-proof and germ killing are poor. Therefore, it is apparently that the wind outlet panel is necessary to be improved so that it is more safer and easily detached and has the function of killing germs."}
{"text": "1. Field of Invention\nThis invention relates generally, to a method for forming a doped region in a semiconductor device, and more particularly, to a method for forming a shallow doped junction using ion implantation.\n2. Description of the Prior Art\nAs the demand for high performance semiconductor devices increases, device manufacturers continually redesign semiconductor devices, such as integrated circuits, to have smaller and smaller dimensions. Of the many different device components used in an integrated circuit, the dimension of the gate electrode in a metal-oxide semiconductor (MOS) device is the benchmark dimension by which all other device components are measured. Typically, state of the art MOS devices have gate dimensions in the range of about 0.1 to 0.25 microns.\nAs gate dimensions are scaled to smaller and smaller values, a concomitant requirement for a reduction in the junction depth of transistor components such as source and drain regions, exists. In addition to requiring a small junction depth, state of the art MOS transistors also require extremely high surface doping concentrations in the source and drain regions. Typically, many doped regions are formed in a semiconductor substrate in the vicinity of the substrate surface. Dopant atoms of either n-type or p-type are typically incorporated into the silicon substrate using both thermal diffusion and ion implantation.\nIon implantation is a physical process in which dopant atoms are ionized, accelerated to a velocity high enough to penetrate the surface of a silicon substrate, focused into a narrow ion beam, and scanned as a beam across the surface of a semiconductor substrate. Dopant ions impacting the surface of the substrate enter the substrate and come to rest below the substrate surface. The depth of penetration into the semiconductor substrate depends upon the particular species and the ion implantation energy. Ion implantation is used in most doping operations in the fabrication of submicron dimension integrated circuit devices. State of the art ion implantation processes can be carried out to form doped regions of precise doping concentration and implantation depth. Doping by means of thermal diffusion is a chemical process typically used to dope thin film layers that do not require the formation of well-defined doped regions.\nTypically, source and drain regions in MOS transistors are formed by ion implantation techniques, and in the case of an n-type transistor, dopants such as phosphorus, arsenic, and antimony are commonly used. Modem ion implantation systems analyze an ion beam that has been extracted from an ion source. The extraction efficiency varies depending upon the particular source material and the particular dopant species to be implanted. For common n-type dopants the maximum beam currents are obtained from singly charged individual ions.\nTo operate an ion implantation system efficiently, and at low cost, a maximum beam current is necessary to maintain a high throughput. Accordingly, the formation of n-type source and drain regions is most commonly carried out using singly-charged, single n-type ions. Since the ion implantation depth, for a given species, varies directly with the ion implantation energy, a very low implantation energy must be used to form a shallow junction in the substrate. The ion implantation system must be operated at very low voltages to obtain the low implantation energy necessary to form a shallow doped region. However, both ion beam extraction efficiency and dose measurement can be adversely affected at low voltage operation. Additionally, at low voltages, beam stability is reduced and spot size control becomes more difficult. Accordingly, a need existed for a method of forming a shallow, highly-doped region in a semiconductor substrate, while maintaining optimal ion implantation operating conditions."}
{"text": "1. Technical Field\nThe present application relates to a measurement system.\n2. Description of the Related Art\nA technique to visualize or a technique to detect a pedestrian and an obstacle (hereinafter may be referred to as an “object”) is important to improve the safety of automobiles. An image pickup system using a range gate camera is known as a representative instance. The range gate camera is a camera that is designed to have exposure (the gate is opened) only for an extremely short specific time. It is possible to pick up an image of only an object at a predetermined distance by a combination of a range gate camera and pulsed irradiation light. With a range gate camera, an adverse effect of a fog which is present in front of an object may be eliminated. A range gate camera also functions as a time of flight (TOF) image pickup system by successively changing gate opening timing. A range gate camera is disclosed in, for instance, in Mitsubishi Heavy Industries Technical Review Vol. 42, No. 5, P. 212 (2005-12)."}
{"text": "The present invention relates to an instrument for the ligation of hemorrhoids or the like.\nVarious methods for the therapy of hemorrhoids are currently used when said hemorrhoids no longer respond to strictly medical cures or when their size becomes such that any non-surgical therapy would produce no result or even negative results.\nA first method is the strictly surgical one, which is irreplaceable in case of considerable internal and external prolapses.\nThis is a painful method which entails significant postoperative hospitalization and is also usually rather destructive as regards the muscles of the anal-proctorectal tract.\nA second method is constituted by cryosurgery, which in practice consists in \"freezing\" the hemorrhoids at temperatures which can vary between -59 degrees Celsius and -89 degrees Celsius for cycles lasting a few minutes.\nAnother method is constituted by thermosurgery, which uses devices which emit flows of heat at preset temperatures and for limited periods of time.\nA method which is currently having considerable success, since it is the easiest, quickest and most practical, as well as the least painful and least damaging, is the so-called \"throttling\" surgery.\nThis method consists in applying an elastic ring of adequate size on the \"neck\" of the hemorrhoid, leaving it in place until the hemorrhoid necrotizes and falls away spontaneously due to the occlusion of the blood supply pathways.\nThe time required in order to achieve this effect is relatively short, the results are extremely positive and the method, which is typically ambulatory, is not painful.\nA cylindrical bush is used to apply the ligature to the \"neck\" of the hemorrhoid, and the elastic ring is arranged around said bush, dilating the ring considerably, by means of a conical applicator.\nOnce the hemorrhoid has been passed through said bush, the elastic ring is caused to slide, moving it so that it surrounds the \"neck\" of the hemorrhoid and, by pressing on it, performs the throttling action.\nTwo methods are currently used to make the hemorrhoids pass in the bush.\nA first one of said methods consists in making a special instrument pass through the bush; said instrument is provided with a forceps at its end, which is suitable for gripping a portion of the hemorrhoid and for pulling it inside the bush.\nThis method has a series of problems, some of which are considerable; first of all, the need arises to use both hands, one for the bush and one for the instrument provided with the forceps.\nIt is furthermore rather traumatic, since the forceps can break the wall of the hemorrhoid.\nFurthermore, it is not always possible to achieve the passage of the entire mass to be throttled.\nIn a second method, which is pneumatic, the bush constitutes the end of a duct which is connected to an aspirator which, when actuated, aspirates the hemorrhoid inside said bush.\nAlthough this method has the advantage that it does not break the wall of the hemorrhoid and can be performed with one hand, it has the problem that during the operation the bush is constantly at negative pressure, so that the hemorrhoid cannot be disengaged in case of incorrect grip; furthermore, at the end of the operation it is necessary to disconnect the aspiration system from the instrument in order to be able to disengage the hemorrhoid.\nFor both \"throttling\" methods, the typically metallic material of which the instruments are made must be resterilized after each use and, since such materials are not transparent, it forces the operator to work practically blind.\nThe instruments themselves are furthermore considerably expensive."}
{"text": "In the editing system, it has been a conventional practice to make various processing for image signals by the key signal which is generated with reference to the luminance signal as a standard.\nSpecifically, the editing system of this type extracts the luminance signal form an image signal which is an object of processing and makes the gain of the luminance signal variable after clamping it to generate the key signal. Further, with reference to the key signal, an image signal of background and the image signal to be processed are mixed each other, thereby causing the image as the object of processing (see FIG. 14B) to be fit into the image of background (see FIG. 14A) for producing a composed image (see FIG. 14C), for example.\nIn the meantime, the editing system may generate the key signal taking an image signal produced by the computer graphics as the object of processing. In such an image signal as produced by the computer graphics, there may sometimes be included a picture of blue color, etc. having little luminance signal components.\nIf the key signal is generated with the luminance signal as a standard in connection with the picture of blue color, etc. having little luminance components, the key signal will be generated in thinner form. Consequently, when processing the picture produced by the computer graphics by the key signal, the edition processed picture will be displayed played in the thinner shape. This will make the edited picture to cause a sense of incongruity.\nThus, when letters of Y and C in FIG. 14B are displayed in white and blue, respectively, there is an inconvenience that, as is shown in FIG. 14C, the letter C will be composed to be thinner than the letter Y in the composed picture."}
{"text": "1. Field of the Invention\nThe instant disclosure relates to a solar greenhouse structure; and more particularly, to a solar greenhouse having a light weight and configurable structure.\n2. Description of Related Art\nWith the increasing environmental awareness, solar power has been widely adopted for practical use. Solar panels in sophisticated modern devices, such as solar batteries or solar heaters, are used to convert solar radiation into useful forms of energy such as electricity or heat. To effectively obtain solar energy, the solar panels need to be mounted at unobstructed locations, usually in an elevated manner, by specially designed supporting structures.\nA greenhouse is commonly used as an agricultural shelter having different types of covering materials that are usually transparent to solar radiation. Due to the nature of this agricultural structure, greenhouses are usually built in locations where sunlight is abundant. Thus, intuitively, the rooftop of a greenhouse structure is an ideal place for mounting solar panels. However, most conventional greenhouses are not designed to host solar panels. Particularly, most light-weight or low-budget greenhouses, lack the structural integrity to host solar panels on their rooftops. As a result, the remodeling cost for converting conventional greenhouses is often very high.\nIn addition, externally mounted solar panels must be rigid enough to withstand adverse weather conditions, such as rain and hails. Therefore, the external mounting of solar panels usually requires the use of thicker solar panels for added strength or additional protective layer. As a result, the cost of setting up a solar greenhouse is relatively high. The attempt to increase structural strength inevitably adds weight to the solar panel.\nHowever, conventional greenhouses often use vinyl coverings; their main structures are usually not designed to withstand excessive weight loading. As a result, the additional weight on the solar panels makes them less suitable for external installation on conventional greenhouse structures, particularly to the roof structure thereof.\nTo resolve the above issues, the inventor researched and proposes the following design solutions to address the above shortcomings."}
{"text": "1. Field of the Invention\nThe present invention generally relates to computer image generation for display to a user and, more particularly, to an architecture for a compositing subsystem for rendering images of opaque, translucent, transparent or semi-transparent items or objects which may be overlaid on the same display space, especially where the screen is shared by multiple tasks in a multi-tasking environment.\n2. Description of the Prior Art\nIt has long been recognized that visual displays are much preferred for a communication interface between a data processing device or system and an operator. Visual displays provide the operator with confirmation of input actions as well as appropriate input options and instructions, often in the form of icons and/or menus and to exert control over plural applications which may be run concurrently on the data processing system, such as in a multi-tasking environment.\nVisual displays have become fairly sophisticated with various display attributes used to enhance communication with a user. Color icons and animation have been used to enhance visibility and indicate both cursor function and identity (e.g. the input device controlling the cursor) as well as system function. The appearance of the visual display, including various display features which may be provided, is a major factor in the \"look and feel\" of particular data processing applications. By the same token, the appearance of the display, particularly as it is changed in response to user commands or system function or operation is a major factor in user confidence and satisfaction with both the data processing system and applications run on it.\nUtilization of multiple display features by one or more applications also calls for an order to be observed in their presentation (often referred to as \"compositing\", at the level of display screen image generation) to achieve appropriate rendering of each image portion. For example, in a windowing display environment, an order must be established between displayed, overlapping windows so that data being manipulated in one active window is not obscured by objects associated with another window. Menus, at least in the active window, should not be obscured. Cursors must not be obscured by the data on which they operate but, except for a so-called main cursor, may be covered by windows or menus. Sprites may be used both as cursors and application level objects, such as icons (with or without animation) with corresponding visibility requirements.\nAdditional complexity in ordering results windows which are often used in multi-tasking environments. Windows, themselves, may be regarded as either opaque or transparent. Further, if graphics applications are run in any window, that window may display many objects in many colors (as well as cursors, menus and other display features) which may overlap (e.g. occult each other in various combinations and orders) and which may also have various degrees of opacity or transparency (e.g. translucency) as well as color, pattern, texture, intensity, animation and other display attributes.\nAnother complexity derives from the expected behavior of display image features. For example, pop-up features such as help windows or balloons and pull-down menus are generally stationary within a window but are expected to appear and disappear from a window almost instantaneously and without leaving artifacts such as blank spaces or obsolete images (e.g. if the image of an object should have been changed while obscured by another images). Images of graphic objects may be dragged from one location to another or be distorted or changed in size by dragging a handle of a border, corner or some other portion of the object and portions of objects or background which have been obscured should immediately reappear, correctly rendered, when the object no longer occults them. Similarly, images of three-dimensional objects may also be rotated in space so that the projection in the plane of the screen may cause the displayed boundary of the object to change. While users do not generally object to the time required for redrawing a screen when an object is added to or deleted from a screen or window, immediate response is expected during many other user interactions with the data processing system.\nDisplay generation is a computationally intensive process in view of the number of pixels to be controlled at high resolution of color, intensity and hue and the number of objects, display features and screens and interactions therebetween which must be supported. Therefore, display generation is usually very expensive in processing time and image data storage. Therefore, buffering of all objects and display features (even though obscured in the final image) does not provide a viable solution to rapid response time of the display. Such potentially massive storage requirements require significant processing time for management and access.\nFurther, full generality of rendering expands computation time exponentially with the number of layers of data which must be considered in developing the correct pixel image values for control of the display device since each underlying layer can potentially affect the rendering of all layers above it. That is, the compositing process must support simulation of the effects of overlaying of semi-transparent colored filters so that a display may be formed of opaque or semi-transparent shapes that may appear to float over and under each other without flickering or impairing rendering of graphics beneath them. To date, even modest levels of generality of compositing have required substantial hardware provision or such amounts of image processing time as to place sophisticated image processing within acceptable time constraints beyond the capacity of all but the largest of data processing systems."}
{"text": "In the office equipment industry, particularly in regard to what are generally called “copiers,” much effort is directed to enabling a human user to communicate to the machine the actions that are desired. Even in the basic copying context, a human user will have to enter such information as what size original documents are being loaded, what size output sheets are desired, whether transparencies are desired, whether the copies are to be one- or two-sided, how many copies are desired, how are the copies to be stapled, are the images to be changed in size, etc. Various schemes have been proposed over the years for presenting possible selections to a human user, usually through a changeable display such as a touchscreen, in a convenient and intuitive way. A patent giving an overview of “user interfaces” in a copier is U.S. Pat. No. 5,159,546.\nIn recent years the basic office copier has evolved into what can be called a “multi-function device” or MFD. With digital technology, a machine with the basic outward appearance of a traditional copier can perform additional functions such as printing documents submitted in digital form over a network; sending and receiving messages via facsimile; recording hard-copy original images in digital form and sending the resulting data over a network, such as in electronic mail; or recording hard-copy original images in digital form on a compact disc or equivalent medium. These additional functions present further challenges to a designer of a convenient and intuitive user interface.\nWhen a human user approaches a digital copier or MFD, there are two basic actions: a document of hard-copy image-bearing sheets is loaded into the machine, and the user interacts with a series of screens at the user interface (UI) to communicate to the MFD the desired actions. Typically, the entry of selections into the UI consumes a significant amount of time, and is prone to mistakes and changes of mind on the part of the human user. Typically, the recording of the images on the original sheets cannot begin until the control system of the machine has been instructed through the UI."}
{"text": "Invasive manipulation of the peripheral circulatory system to correct occlusive diseases of the arteries has become more and more routine. Over 200,000 procedures are performed annually in the United States alone which involve blood vessel bypasses, balloon catheters, and other \"mechanical\" techniques to correct the problem. A serious side effect of these procedures is the subsequent development in the subject of intimal hyperplasia which may, itself, constitute a blockage problem. This appears to be a direct response to the intimal injury caused by the intervention; smooth muscle cells and fibroblasts proliferate and create stenoses in the interior of the vascular wall. The mechanism of this process is evidently not well understood, but a central feature of the development of the problem appears to be unwanted proliferation of smooth muscle cells.\nVarious methods have been tried to overcome this dangerous side effect. Approaches have included mechanical manipulation as well as administration of chemical agents such as aspirin, dexamethasone, heparins, calcium channel blockers, and a variety of other agents presumed on the basis of various theories to interfere with the development of the stenoses. They have met with very little success and have side effects of their own.\nOne additional approach is the use of photodynamic therapy (PDT). This form of management, originally applied to cancer treatment, involves the use of photoactive materials which home to tumor tissue, presumably because of the rapidly proliferating nature of the tissue. The photoactive substances, which include psoralen, various porphyrin-based materials, such as Photofrin II.TM. porphyrin aggregate, chlorins, phthalocyanins, and monohydrobenzoporphyrin derivatives, to name but a few, are harmless unless photoactivated. However, when irradiated with light of appropriate wavelength, the drugs apparently effect the formation of a toxic agent, presumably singlet oxygen, although they themselves are chemically unchanged. The resultant toxic agent causes the destruction of the unwanted tumor tissue.\nPDT has also been applied with some success to the treatment of atherosclerotic plaques. See, for example, Kessel, D. et al., Photochem Photobiol (1984) 40:59-62; Okunaka, T. et al., Photochem Photobiol (1987) 46:769-775; Spears, J. R. et al., J Clin Invest (1983) 71:395-399; Spokqiny, A. M. et al., J Am Col Cardiol (1986) 8:1387-1392; Copperath, K. et al., Eur Heart J (1989) 10 (Suppl):151; Straight, R. et al., Photodynamic Therapy of Tumors and Other Diseases, (1985) pp. 349-350.\nThe use of PDT to treat or prevent the restenosis that often accompanies angioplasty has also been studied. These studies have either employed smooth muscle cells in (SMC) in culture, on the theory that SMC proliferation is the sine qua non of restenosis, or have used animal models. In general, PDT appears to show promise in this regard. However, controls run in many of these studies, using the photoactivating agent in the absence of light, have provided contradictory results when smooth muscle cells in culture were used as the model system. Applicants are unaware of any animal studies which showed any indication of positive results for preventing restenosis in the absence of light.\nDartsch, P. C. et al. reported in Advances in Laser Medicine 4: Laser Angioplasty II Biamino, G., et al. (ads) Ecomed Verlagsgesellschaft, Landsberg/Lech, Berlin (1990) 77-80, the results of contracting smooth muscle cells in culture exposed to dihematoporphyrin-ester and -ether (DHE). This porphyrin-based drug is now marketed as Photofrin II.TM.. The report discloses that SMC were isolated by enzymatic disaggregation of either normal or stenosing plaque tissues and cultured in vitro. They were tested in their first, second or third passage by treating them with DHE at concentrations ranging from 0.1-25 .mu.g/ml and irradiated with ultraviolet light. The percentage of viable and still adherent cells was markedly reduced for plaque-derived SMC and much less dramatically reduced for normal SMC. A less dramatic, but nevertheless detectable, effect was observed in the presence of DHE but in the absence of radiation. Specifically, cells treated with 5 .mu.g/ml DHE and light showed a reduction in the number of viable cells to 73% in the case of normal derived cells and 38% for plaque-derived SMC. Using a DHE concentration of 1 .mu.g/ml and an energy density of 1200 mJ/cm2, after 24 hours 80% of the normal SMC and 20% of the plaque-derived SMC were viable. Similar results were reported by this group in related publications: Dartsch, P. C. et al., Atherosclerosis (1990) 10:616-624; Dartsch, P. C. et al., J Am Coll Cardiol (1990) 15:1545-1550.\nIn two recent reports by Sobeh, M. S. et al., results apparently contradictory to those of Dartsch, et al. were obtained when SMC cultured from the intermedia of human long saphenous vein harvested for coronary artery vein bypass grafting, were used in the tests. These cells were treated with Photofrin II.TM. porphyrin aggregates and irradiated. These reports state that the cells were unaffected by Photofrin II.TM. porphyrin aggregate at 0-100 .mu.g/ml without light. However, when treated with light energy of greater than 3 J/cm2 in the presence at least 2 .mu.g/ml of the drug, a mean cell destruction of over 80% was reported regardless of wavelength. The reports also state that light without prior chromophore sensitization produced no cell damage. (Vascular Surgical Society of Great Britain and Ireland Annual General Meeting, London, November 1992; 13th Annual American Society of Laser Medicine and Surgery, New Orleans, April 1993.)\nAsahara, T. et al. reported in Circulation (1992) 86 (Suppl) 1-846, that PDT was able to inhibit restenosis in rabbits that had received balloon injuries of the iliac artery and were fed with a 0.2% cholesterol diet. Hematoporphyrin derivative (HPD) was administered 24 hours before irradiation at various times relative to the injury. The best effects were observed when the treatment was administered one week after angioplasty.\nIn another in vivo study, Ortu, P. et al., Circulation (1992) 85:1189-1196, reported that photodynamic therapy, with chloraluminum-sulfonated phthalocyanine (CASPc) used as the drug, was effective in inhibiting the intimal hyperplasia response in rats subjected to balloon injury of the carotid artery. Controls consisted of rats irradiated, but not administered the drug. No controls using CASPc without light were reported.\nEton, D. et al., J. Surg. Res. (1992) 53:558-562, reported the effect of photodynamic therapy using Photofrin II.TM. porphyrin aggregates in a rabbit model wherein the rabbits underwent standardized intimal injury to both common carotid arteries with a balloon catheter. The test animals received Photofrin II.TM. porphyrin aggregate and subsequent irradiation; the control groups either received no treatment, or chromophore alone, or light alone. The results were evaluated in terms of arterial cross sections. Only the test group showed a statistically significant improvement over the controls, although the animals treated with light alone or Photofrin II.TM. porphyrin aggregate alone had non-significant lower mean ratios of the area of intimal hyperplasia to the area enclosed by the internal elastic lamina, used as a measure of stenosis in this study.\nIt would be advantageous to provide a treatment to prevent intimal hyperplasia (IH) following vascular trauma which is independent of PDT, so that the necessity to provide light using specialized equipment is avoided. It has now been found that green porphyrins, administered concurrently with, and for a period of time after, angioplastic procedures, can effectively inhibit the undesired stenoses often accompanying this procedure."}
{"text": "The invention relates to miter saws and in particular to a miter saw and a fence assembly for supporting a workpiece during cutting of the workpiece.\nCompound miter saws, typically include a base having a support surface for supporting a workpiece and a turntable mounted on the base for pivotal movement about a vertical axis. A saw arm is supported by the turntable for movement therewith, and a fence is supported by the base to provide support for a workpiece on the support surface during cutting of the workpiece. The saw arm has a saw blade and is pivotable about a horizontal axis between non-cutting and cutting positions. If it is desired to engage in compound miter cutting of the workpiece, the saw arm is pivoted about a horizontal bevel angle axis to allow the user to make the angular cuts in the workpiece. In some prior art miter saws, to provide sufficient clearance for the blade of the saw arm and yet provide adequate vertical support to the workpiece during the cutting of the workpiece, the end of the fence adjacent the saw blade has an angular edge portion. The angular edge portion is intended to provide support for the workpiece while also providing clearance for the saw blade during compound miter cutting of the workpiece. It is known in the art to adjust the position of the fence toward and away from the saw blade to allow cutting at extreme bevel angles and yet provide support to the workpiece."}
{"text": "The present invention relates to a semiconductor memory comprising techniques that may be adapted advantageously to an image memory incorporating an operation circuit.\nIn the related art, U.S. Pat. No. 4,766,570 issued on Aug. 23, 1988 discloses an image processing memory that incorporates an operation function and a serial output function."}
{"text": "Rotary lawn mowers and grass cutting blades serving to cut grass or other vegetation are well known in the art today. Typically a grass cutting blade, also referred to as a \"cutting bar\" is rotated about a generally vertical axis within a downwardly opening housing in order to cut grass by rapid vertical movement of one or more cutting edges defined by the cutting blade. The rotating blade may also include one or more air-fanning \"foils\" or \"vanes\" which serve to create air movement into, about, and out of the cutting housing to transport cut clippings either out of a discharge opening (where it can be bagged or broadcast sidewardly away from the housing), or downwardly onto and into the cut grass, where it preferably may be allowed to decompose.\nIn the instance where the grass clippings (or other cut organic matter such as leaves, weeds, etc.,) are caused to be transported downwardly onto the cut grass, it has been found to be desirable to cut the clippings as finely as possible, and to dispense such grass clippings between the cut grass blades and onto the ground surface. This process, often referenced in the art by descriptive terms such as \"mulching\", \"composting\", or \"recyling\", allows the clippings to decompose with minimal interference with the growing of the cut grass. Upon decomposition, the clippings may nourish the growing grass.\nIn the prior art, it has been known to provide different cutting blades and blade housing configurations to adapt to the particular type of cutting desired. For example, certain blades have been found to be acceptable for bagging or side-discharging of clippings, while other blades have been found to be best suited for mulching. However, it is obviously advantageous to provide a single type of blade and accompanying housing which may provide both functions.\nAlthough many different cutting blade configurations are known in the prior art, disadvantages do exist. Some such disadvantages include clogging of the cut grass within the housing or on the blade, due to inadequate cutting of the grass, or inadequate air flow within the blade housing. Other disadvantages include uneven dispersal of the clippings, unacceptable power consumption, and difficulties in resharpening.\nU.S. Pat. No. 3,780,509 to Woelffer entitled \"Rotary Lawn Mower Cutter Bar Width Bi-Level Cutter\", issued Dec. 25, 1973, discloses a grass cutting blade including a pair of cutters on each opposite end of the blade. The pair of cutters includes an upper and lower blade, with the upper blade in rotational advance of the other blade, so that the leading upper cutter will cut the grass prior to being cut by the trailing lower cutter. A air-fanning blade is also provided to create a flow air to convey the clippings from the blade and through the outlet of the mower.\nU.S. Pat. No. 2,760,327 to Bovee, entitled \"Housing for Disc Type Power Mower\", issued Aug. 28, 1956, discloses a lawn mower which includes a pair of cutting blades 31, 30, the upper of which may be in slight rotational advance of the lower blade 30.\nU.S. Pat. No. 3,477,214 to Rogers, entitled \"Rotor Assembly for a Shredder Implement\", issued Nov. 11, 1960, discloses a shredder implement including multiple cutting blades, some of which are vertical spaced from the other and may pivot about their own axis.\nU.S. Pat. No. 2,547,540 to Roberts, entitled \"Power Mower\", issued April 3, 1951, discloses a power mower including a cutter bar 33, itself including multiple cutting blades 36.\nU.S. Pat. No. 2,697,322 to Watrous, entitled \"Cutter Member for Rotating Disk Type Lawn Mowers\", issued Dec. 21, 1954, discloses a power lawn mower including multiple cutting blades such as 15, 16, which may be vertically spaced apart. The cutting blades 15, 16, include corresponding air propelling flanges 18 at their rear margins.\nU.S. Pat. No. 2,908,128 to Mauro, entitled \"Cutting Device for Rotary Mowers\", issued October 13, 1959, discloses a rotary mower including multiple cutting edges.\nU.S. Pat. No. 3,703,071 to Anderson, entitled \"Rotary Blade for Mulching Lawn Mower\", discloses a rotary blade for a power-driven lawn mower which includes an elongate body terminating at each end in a plurality of integral cutting fingers parallel to and displaced successively below and above the plane of the body of the rotary blade.\nU.S. Pat. No. 4,318,268 to Szymanis, entitled \"Mulching Lawn Mower\", discloses a mulching lawn mower which includes a toroidal portion bounded by concentric skirt walls wherein the cross sectional area of the toroid increase towards the inner skirt. A cutting blade 26 is twisted in a particular configuration to provide flipping movement as shown in FIG. 2.\nU.S. Pat. No. 4,977,735 to Davis, entitled \"Lawn Sculpturing Method and Apparatus\", issued Dec. 18, 1990, discloses several grass cutting blades which may be used to provide sculpturing effects in lawns, due to multiple cutting edges at various vertical heights.\nU.S. Pat. No. 2,682,142 to Clark, entitled \"Cutter Element for Rotary Mowers\", issued Jun. 29, 1954, discloses cutting elements in spaced-apart relationships to provide an improved cutting path. FIG. 4 illustrates the interaction of two such cutting blades, one having one cutting edge at each end, and the other having two cutting edges with the upper cutting edge being provide to reduce \"streaking\", or unmown areas between the adjacent cutter units during turning operations.\nU.S. Pat. No. 3,080,697 to Mauro, entitled \"Cutting Blade for Rotary Mower\", issued Mar. 12, 1963, discloses a lawn mower cutting blade which includes a plurality of cutting edges, such cutting edges in spaced-apart relationships to provide different cutting actions depending on the rotation of the reversible cutting blade.\nU.S. Pat. No. 2,857,729 to Zoldok, entitled \"Rotary-Type Mower Blade\", issued Oct. 28, 1958, discloses a mower blade including not only vertical cutting edges, but also inclined cutting edges.\nU.S. Pat. No. 3,538,692 to Cope, entitled \"Grass Mowing and Mulching Device\", issued Nov. 10, 1970, discloses a grass cutting blade including a plurality of cutting edges. The blade includes not only at each end a horizontal cutting edge 4, as well as a plurality of chopping edges 7, positioned behind the cutting edge 4.\nU.S. Pat. No. 4,628,672 to Jones, entitled \"Rotary Cutter\", issued Dec. 16, 1986, discloses a cutter blade including multiple cutting edges, including slicing edges at the extreme ends of the blade.\nU.S. Pat. No. 2,669,084 to Warren, entitled \"Rotary Mower Blade\", issued Feb. 16, 1954, discloses a cutting blade including vertical pivotable cutting blades 15, each of which include a trailing curved portion 23.\nU.S. Pat. No. 2,891,369 to Rietz, entitled \"Rotary Shredder and Cutter\", issued Jun. 23, 1959, discloses a farm-type implement including a shredder bar member, which includes a downwardly-turned hook 215, sharpened at its leading edge as shown in FIG. 5.\nU.S. Pat. No. 3,220,170 to Smith, entitled \"Vegetation Mower and Blade\", issued Nov. 30, 1975, discloses a cutting blade including cutting edges and an upwardly-turned flanges to provide air flow.\nReissue Patent No. 32,377 to Seyerle, entitled \"Cutter Blade for Pneumatically Transporting Grass Clippings\", reissued on Mar. 24, 1987, discloses a cutter blade including an upturned flange for transporting grass clippings.\nAs may be seen, it is known to provide various cutting blades or cutter bars which include multiple cutting edges and which may include upturned or down turned air foils to provide improved air flow. However, disadvantages do exist in the prior art.\nTherefore it may be seen that a need exists for a grass cutting blade which may be used to bag, broadcast, or mulch grass clippings, which cuts the clippings finely. A need also exists for a grass cutting blade which eliminates or minimizes clogging of the clippings on the blade or within the housing, and which provides uniform dispersion of the clippings below the cut grass surface when in its \"mulching\" mode. A need also exists for a grass cutting blade which is easily made by mass production techniques, and which may be easily resharpened. Finally, a need exists for a grass cutting blade which is efficient in its use of power under a variety of cutting conditions."}
{"text": "A standard harvesting machine such as a field chopper has a pickup head that cuts a standing crop or picks up a windrow of previously cut crop, for instance corn, wheat, or soy, and a conveyor mechanism that feeds a stream of the cut crop back through treatment devices that typically crush and comminute it. The conveyor mechanism normally includes at least two rollers that are oppositely driven.\nIn order to protect the machinery from hard objects such as stones that might be entrained in the crop stream, one of the rollers is normally mounted so that it can be deflected against a spring force away from the other roller. Thus when a hard object is drawn in, the one roller will deflect laterally and allow it to pass without damage to the conveyor rollers.\nWhile such a system protects the conveyor rollers themselves from harm, it does not protect the downstream equipment which is equally likely to be damaged by contact with the hard object. Therefore systems are provided that allow the equipment to be stopped or shut down when a hard object is encountered by the intake conveyor.\nIn East German patent 110,413 and 116,553 of Wehsely mechanical feelers are provided between the feed rollers that detect hard objects and operate mechanism for sorting out any hard object. Such systems require a relatively slow feed rate that greatly reduces throughput of the machine. In addition they are expensive and failure prone.\nAnother solution proposed in East German patent 114,893 of Wehsely proposes directing a bright light through the crop stream upstream of the feed rollers, since the crop is normally somewhat translucent and most hard objects are solidly opaque. This system requires that the light be very bright, so that when used with a dry crop like straw there is a severe danger of fire, making the system unsuitable for much field use.\nIn East German patent 111,534 also of Wehsely the pickup head's transverse conveyor is provided with a plurality of electronic sensors that if physically contacted by a hard object emit a signal that is used to stop the system and even eject the object. The problem with this system is that the foreign object is often masked in the crop stream so that it does not contact the sensors which, instead, remain in contact with the relatively soft crop stream.\nU.S. Pat. No. 4,353,199 of Chow describes a harvesting-machine stone detector wherein the feed rollers are provided on their inner surfaces with piezoelectric acoustic sensors that produce a characteristic signal when the roller is struck by a hard object such as a stone. The resultant electrical signal is processed so that when a stone is detected the machine can be shut down or the crop stream can be temporarily diverted out of the machine to eject the stone. The orbiting piezoelectric sensors must be connected through failure-prone commutator rings with the processing circuit so that there is frequent failure of transmission of the signal to the circuit. Furthermore when the stone is well cushioned in the crop stream the sensors will not detect it.\nGerman utility model 296 16 473 of Weiss describes a system where one of the intake rollers can be deflected transversely to the crop stream against the force of a fluid-type shock absorber. A pressure sensor in the absorber can detect the sudden change in pressure caused when a stone wedges between the rollers and forces them abruptly apart. This pressure sensor produces an output that can be processed to trigger stopping of the machine or ejection of the stone when the output of the sensor exceeds a predetermined threshold, indicating that a stone is trying to get between the rollers. The problem with this system is that it is only applicable to high-end machines with pneumatically or hydraulically damped feed rollers. Furthermore it is possible for a sudden pressure peak to be created by something other than a potentially dangerous stone, so that the system gives frequent false alarms.\nFurther such systems are described in European patent publications 0,217,417 and 0,217,418 both of Weiss which measure the acceleration of the feed rollers. The signal is processed, however, identically as in the above cited utility model, producing the disadvantage that, when the crop stream is moving slowly for fine chopping of the crop, a stone can get past the rollers while when the crop is moving rapidly there are frequent false alarms."}
{"text": "The present disclosure relates to stretch leather composites, and more particularly to methods of making stretch leather composites including an animal skin laminated with a stretch fabric.\nStretch leather composites including a soft leather laminated with glue to an elastic fabric are used in various consumer products, such as footwear, apparel, clothing and the like, as an alternative to other stretch materials having less desirable aesthetic characteristics. Existing methods for preparing stretch leather composites include laminating an elastic fabric to leather with glue. These methods include inefficiencies and processing challenges introduced at least by variations in degree of stretch of the elastic fabric during the lamination process and difficulties inherent in handling the stretch fabric during the lamination process. Moreover, stretch leather composites produced by pre-existing methods can lose their shape after being repeatedly stretched.\nAccordingly, there remains room for improvement to produce stretch leather composites that address these disadvantages of pre-existing methods."}
{"text": "This application relates to resonators for electromagnetic signals such as optical resonators and microwave resonators.\nAn optical material such as a dielectric material may be shaped to construct an optical whispering-gallery-mode (“WGM”) resonator which supports a special set of resonator modes known as whispering gallery (“WG”) modes. These WG modes represent optical fields confined in an interior region close to the surface of the resonator due to the total internal reflection at the boundary. Microspheres with diameters from few tens of microns to several hundreds of microns have been used to form compact optical WGM resonators. Such spherical resonators include at least a portion of the sphere that comprises the equator of the sphere. The resonator dimension is generally much larger than the wavelength of light so that the optical loss due to the finite curvature of the resonators is small. As a result, a high quality factor, Q, e.g., greater than 109, may be achieved in such resonators. Hence, optical energy, once coupled into a whispering gallery mode, can circulate within the WGM resonator with a long photon life time. Such hi-Q WGM resonators may be used in many optical applications, including optical filtering, optical delay, optical sensing, lasers, and opto-electronic oscillators."}
{"text": "1. Field of the Invention\nThe present invention relates to a numerical simulation in a real time system and particularly to a high speed processing apparatus which is suitable for computing a system having sparse and irregular characteristics such as an electronic circuit, a logic circuit, a structure analysis, a tidal current computation or the like.\n2. Description of the Prior Art\nAs a computer for rapidly performing a numerical computation, there exists super computers and array processors, however they have the following problems.\n(1) In case of accomplishing a computation in a sparse or irregular system as described above, the efficiency is remarkably reduced because of shortened vector length. PA1 (2) It is difficult to increase the degree of parallel processing so as to speed up the computation. PA1 (i) The kinds of nodes are limited to the four rules of arithmetic. PA1 (ii) There exist no loops from a view of a graph. PA1 (iii) There is no case of converting a value of the input terminal (a.sub.11, a.sub.31, . . . in FIG. 1(a)). PA1 (i) The input data are assembled, that is, the operation-executed nodes are discriminated. PA1 (ii) The operations of these first nodes (only the four rules of arithmetic) are executed. PA1 (iii) The nodes to which the outputs of the first nodes are searched and the processing results of the first node are set as the input data of the nodes to which the results of the first nodes are output.\nThough there are some realized devices in a data flow machine (Mizoguchi, et al.: Image processing processor TIP-1, Shingakugiho: EC81-59), there is a limit to the speedup because the data reading section is series-connected."}
{"text": "The present invention relates to an image processing device and a system using the same, which is capable of performing graphics processing, drawing processing and display processing at high speed.\nMore specifically, the present invention relates to an information terminal machine, such as a facsimile machine, printer, graphics device, portable personal terminal machine, navigation machine and amusement device, and to an image processing system which performs inputting, processing, editing, accumulating, communicating, outputting and displaying of image data using the information terminal machine.\nIn particular, the present invention relates to an image processing device and a system using the same which can minimize the amount of data being transferred between a memory device and the image processing device to thereby perform high speed processing.\nFurther, in particular, the present invention relates to a graphics drawing method, and to an image processing device and a system using the same in which, in order to minimize the amount of data being transferred between a memory device and the image processing device, reduces drawing suppression processing (hidden surface processing) due to graphics overlapping, and which, in particular, performs complex image processing and three dimensional graphics display processing at high speed.\nFurther, in particular, the present invention relates to a method which suppresses an increase of components in an image processing system having high speed access to a memory, thereby reducing the construction cost thereof.\nVarious conventional systems for effecting data transfer between a memory, which holds image data, and an image data processor and/or an image processing device, which processes the image data, have been developed. However, these conventional systems have inherent problems with regards to their processing speed, such as for accessing and drawing, in that high speed processing and real time processing in a device or a system which is required to process a great amount of data have proven to be insufficient.\nA proposal in which image data transferred between a memory device and an image processing device is drawn and processed in blocks representing a pixel aggregate is disclosed in Andy Goris et al. xe2x80x9cA Configurable Pixel Cache for Fast Image Generationxe2x80x9d (IEEE, Computer Graphics and Application, May 1987, pp. 24-32), which is hereinafter referred to as the Goris et al. publication.\nAccording to Goris et al., pixel data is fetched in blocks on demand into a pixel cache, and a drawing process is executed such that a prefetch effect is limited within a block unit.\nEven in a case when a drawing processing is performed while crossing the boundary between blocks, since the pixel data is generated via a common rasterizing mode, a prefetch for pixel data in the subsequent block is started.\nTherefore, if the capacity of the pixel cache is small, in response to the prefetch of the pixel data in the subsequent block, the pixel data in the previous block is first pushed out of the pixel cache and then the pushed out pixel data again needs to be accessed via rasterization, which causes a problem in that frequent memory access is necessitated.\nIn particular, when a coloring-up darkening process is performed for a triangle, although it is necessary to generate pixel data two-dimensionally and to process the same, the memory access is performed one-dimensionally; therefore, even when performing a drawing processing of nearby pixel data, the prefetched data as indicated above cannot be utilized, which causes an inefficient memory access. In other words, it causes a problem in that an additional memory band width is necessitated.\nFurther, in a graphics architecture in which a memory area for drawing processing is arranged in a main memory, a sufficient memory band width has to be provided which causes a problem in that the memory can not be efficiently used.\nA three-dimensional graphics display device generally uses a display method in which an object to be displayed is divided into small three-dimensional graphics, such as triangles and quadrangles, and is modeled; a geometrico arithmetic operation is performed on these small graphics based on their viewed direction; and, based on the resultant arithmetic values, the respective graphics are projected on the two dimensional coordinate of a display device, such as a CRT. At this time, however, it is required to check for overlapping of the respective graphics in their depth direction and to prevent drawing of hidden graphics on a frame buffer.\nAs a conventional method of judging graphics, an overlapping Z buffer algorithm, in other words a depth buffer algorithm, is generally used. This conventional method is, for example, explained in W. M. Newman et al. xe2x80x9cDialog Type Computer Graphics (II) 2nd editionxe2x80x9d translated by Setsuo Ohsuga (McGraw-Hill Book Company, 1984, pp. 483-441), which is hereinafter referred to as the Newman et al. publication. According to this method, depth information is provided for respective pixels every time. When drawing respective pixels, depth information of a pixel to be drawn is compared with the depth information of already drawn pixels, and when the pixel to be drawn is located closer to the screen, the pixel is permitted to be drawn, and conversely, when the pixel to be drawn is located deeper than the already drawn pixel, the drawing of the pixel is prevented.\nAnother method in which overlapping of graphics is checked via geometric calculation is, for example, explained on pages 442-443 of the Newman et al. publication. In this method, a circumscribed rectangle is defined for each of the polygons to be drawn and overlapping of these circumscribed rectangles is judged; however, when the respective circumscribed rectangles do not overlap each other, the judgment with regard to their depth is omitted.\nOn the other hand, a method which avoids the necessity of the Z buffer is explained in James D. Foley et al. xe2x80x9cFundamentals of Interactive Computer Graphicsxe2x80x9d translated by Atsumi Imamiya (Published by Japan Computer Association, 1982, pp. 569-572), which is hereinafter referred to as the Foley et al. publication. This method is based on a depth sort algorithm in which the drawing is started from the graphic in the deepest location and graphics close to the viewer are successively drawn while overlapping graphics.\nFurther, U.S. Pat. No. 4,303,986 discloses a provision of a selective writing means for a memory which stores two dimensional images.\nIn the method disclosed in the Newman et al. publication, a Z buffer algorithm is used for judgment of graphics overlapping, so that this method necessitates depth information for every pixel and therefore requires a Z buffer (frame buffer) of large capacity. For example, assuming a standard system having 34 bits for color information and 16 bits for the Z buffer with a screen having 1024xc3x97768 pixels, a memory of about 4M bits in total is necessitated, among which a memory of about 1.4M bits is needed for the Z buffer algorithm and a memory of about 1.6M bits is needed for the Z buffer. Further, in the conventional Z buffer algorithm, the same pixel is accessed as many times as the number of graphics that overlap, and every time a drawing is performed, the Z buffer content has to be read and the read Z value has to be compared, among which comparison only one comparison result is significant. For this reason, useless memory access is necessitated, which makes an improvement of drawing performance difficult.\nOn the one hand, the method disclosed in the Foley et al. publication, which avoids any need for the Z buffer, is effective with regard to memory capacity reduction; however, since graphics are successively drawn in an overlapping manner from the graphic in the deepest location, the graphic closest to the viewer sometimes can not be drawn depending on the number of graphics to be drawn. Thus, this method can not be applied to a system which requires a real time performance.\nWith regard to a conventional image processing system, for example, JP-A-5-258040 (1993) and JP-A-5-120114 (1993) disclose examples of a data processing system which makes use of a synchronous DRAM permitting high speed data transfer.\nHowever, these image processing systems are not practical because they require excessive time when a plurality of image processings are performed via multi tasks. Further, when a bus control is performed via a time slot method by making use of a synchronous DRAM, a mishit control is frequently generated which causes a problem in that an increase of the throughput is limited.\nFurther, GAIN (Technical Report Published by Hitachi, Ltd. Semiconductor Division, No. 96, 1993.1, pp. 6-11) discloses an example for reducing the size and cost of the device by making use of a built-in RISC (Reduced Instruction Set Computer). However, the paper is silent with regard to a specific memory access method and bus utilization method which realize a high speed image processing.\nStill further, JP-A-4-107056 (1992) discloses a high speed processing method in which a bus which transfers image data from a decoder to a printer is made independent from a MPU bus.\nStill further, in a conventional facsimile machine, an example of an image processing system is disclosed in Shuichi Fujikura et al. A Development of a LSI for Facsimile Image Processingxe2x80x9d (Oki Denki Research and Development Report, October, 1992. No. 156, vol. 59, No. 4, pp. 65-70) having a processor and high speed memory dedicated for image processing in an image input and output unit, and in which data distortion is corrected to thereby realize images of high quality. However, in association with gathering of control units, each formed by a one chip microcomputer, the space rate occupied by the image processing unit increases, which adversely affects the cost thereof.\nStill further, the image processing system in a recent business use facsimile machine tends to accelerate high image quality, high processing speed and large memory capacity such that LSIs dedicated for image processing and for coding are frequently constituted to have their own respective SRAMS. Accordingly, there arises a problem of increased cost of the devices.\nStill further, a conventional image processing system used for a facsimile machine, a printer and a graphics device, as disclosed in JP-A-61-261961 (1986), has a SRAM (static memory) used for local processing by referring to nearby pixels at high speed and a DRAM (dynamic memory) used for storing data, such as symbol data and font data, operating at low speed, but having a large memory capacity. Therefore, the impossibility of integrating the above two types of memories is a significant problem from the point of view of device size reduction, integration into a single LSI, device constitution, device cost and product series development.\nStill further, one of the reasons why high speed image processing could not be achieved with the above conventional art is that the image inputting and outputting processing and the communication processing function are required to have an extremely high real time property as well as a high speed bus throughput of about 4-20 MB/s such that their processings have to be performed via a dedicated processor and a local processing use dedicated memory independent from a main memory.\nDue to the development of the semiconductor micro-machining technology and improvements in microprocessor architecture, a high speed processor, such as a RISC, and a device having an operating speed more than 100 MHz, such as a RAMBUS and synchronous DRAM, have appeared on the market. For example, the synchronous DRAM has already begun to draw attention as a memory having a large capacity operable at high speed. In contrast to the conventional DRAM, a synchronous DRAM can input and output data, address and control signals in synchronism with clocks to thereby realize a high speed data transfer comparable with the conventional SRAM, in addition to the fact that a synchronous DRAM having a larger memory capacity than the conventional DRAM can be realized at a low cost.\nThere are marked characteristics in image communication and processing, including advantageous characteristics with regard to the construction thereof which involve regularity in address renewal, such as the continuity of the address to be processed, easy previous forecasting of processing quantity and simple processing content and limited nearby influence of the processed result, and disadvantageous characteristics representing an intense real time requirement and a possible system break-down when not completing the processing within a predetermined time. However, no devices and systems have been proposed until now which optimize the processing in view of the above advantageous and disadvantageous characteristics, so that it is necessary to provide a device and a system therefor which take into consideration the above advantageous and disadvantageous characteristics.\nAn object of the present invention is to provide a device which processes a great quantity of data at high speed and in real time, and a system using the device.\nAnother object of the present invention is to provide an image processing device which minimizes the amount of data being transferred between a memory device and the image processing device and performs high speed processing, and a system using the device.\nStill another object of the present invention is to provide a graphics drawing method which, in order to minimize the quantity of data transfer between a memory device and an image processing device, reduces drawing suppression processing (hidden surface processing) due to overlapping of graphics and, in particular, performs complex image processing and three dimensional graphics display processing at high speed, and an image processing device and a system using the method.\nA further object of the present invention is to provide a device and a system which avoid an increase of components in an image processing system which accesses a memory at high speed so as to achieve low cost.\nA still further object of the present invention is to provide a graphics drawing method and device, and a three dimensional graphics display method and device which make it possible to reduce the capacity of a data storage area in which drawing use graphics data is to be stored.\nA still further object of the present invention is to provide a graphics drawing method and device, and a three dimensional graphics display method and device which permits high speed drawing processing.\nA still further object of the present invention is to provide a high speed image processing system of low cost by integrating an image inputting and outputting process, a coding and decoding process and a communication process, which are conventionally processed independently, and further by integrating a high speed local processing use memory and a large capacity memory.\nAn aspect of the invention is characterized by the fact that, based on the processing status of pixel data of every one of respective rows in a block buffer which holds pixel data accessed from a memory by a block unit, a command with regard to image processing is executed and the processing of the pixel data is performed.\nFurther, an aspect of the present invention is characterized by the fact that, based on the processing status information of the pixel data to be processed, the pixel data to be processed next is specified.\nFurther, an aspect of the present invention is characterized by the fact that the device according to the present invention comprises a command buffer which holds commands indicating processings relating to drawings, a block buffer which accesses pixel data held in a memory in every block unit and holds the same, a pixel data processing unit which executes the processing of the pixel data in the block buffer based on the command and detects the processing status of the pixel data at boundaries of the respective rows in the block buffer and a status register which holds processing status information indicating the processing status of the pixel data at the boundaries of the respective rows and, in that the processing of the pixel data in each of the blocks is executed based on the commands and the processing status information.\nAn aspect of the present invention is characterized by the fact that, since the commands relating to the image processing are executed and the processing of the pixel data is performed on the basis of the processing status of the pixel data in the block buffer, all of the pixel data in a block unit can be processed without recourse to the order of the raster development, the number of access times between a memory and an image processing processor can be reduced and a high speed image processing and drawing processing can be realized.\nAn aspect of the present invention is characterized by the fact that, since the pixel data to be processed next is specified on the basis of the processing status information of the pixel data to be processed, all of the pixel data in a block unit can be processed without recourse to the order of the raster development, the number of access times between a memory and an image processing processor can be reduced and a high speed image processing and drawing processing can be realized.\nAn aspect of the present invention is characterized by the fact that, since the processing status information of the pixel data to be processed is held in every row unit within each of the blocks, the capacity of the status register can be reduced and the pixel data to be processed next can be specified.\nAn aspect of the present invention is characterized by the fact that, since the processing status information of the pixel data to be processed is held in every pixel data at the boundary with an adjacent block, the capacity of the status register can be reduced and the pixel data to be processed next can be specified.\nAn aspect of the present invention is characterized by the fact that, in order to reduce information for judging overlapping of graphics, a control data writing region is provided indicating whether or not a pixel is already drawn in a data storage region storing drawing use graphics data, control data flags are stored in the control data writing region, and, when the information of the control data writing region indicates that the drawing is not completed, the drawing is performed and otherwise the drawing is prohibited.\nAn aspect of the present invention is characterized by the fact that, when drawing successively a plurality of pixels based on the data of pixel groups relating to a three dimensional image, a control data writing region is provided indicating whether or not a pixel is already drawn in a data storage region storing drawing use graphics data of the respective pixels, and, when drawing of the respective pixels is commanded, a pixel to be newly drawn is designated depending on the control data in the control data writing region, the designated pixels are successively drawn and control data indicating that the drawing is completed is written in the control data writing region relating to the pixel of which drawing is completed.\nAn aspect of the present invention is characterized by the fact that, when drawing successively a plurality of pixels based on the data of pixel groups relating to a three dimensional image, a control data writing region is provided indicating whether or not a pixel is already drawn in a data storage region storing drawing use graphics data of the respective pixels, when drawing of the respective pixels is commanded, and only when control data indicating that a pixel is not yet drawn is added in the drawing use graphics data of the respective pixels extracted from the data storage region, will the drawing of the pixel be executed and control data indicating that the drawing is completed be added in the drawing use graphics data for the pixel and stored in the data storage region; and, when control data indicating that the pixel is already drawn is added in the drawing use graphics data of the respective pixels extracted from the data storage region, redrawing of the pixel is prohibited.\nAn aspect of the present invention is characterized by the fact that the device according to the present invention is applicable to devices in which, based on the data of pixel groups relating to images of a plurality of polygons, the pixels of the respective polygons are successively drawn.\nAn aspect of the present invention is characterized by the fact that the device according to the present invention is applicable to devices in which a graphics group in a form of graphics list containing image information of pixel groups relating to a plurality of polygon graphics having different depths is sorted according to the depth information relating to the depth of the respective graphics and the pixels of the respective graphics are successively drawn according to the order of the sorting.\nFurther, an aspect of the present invention uses a method in which, when performing drawing based on a graphics list containing image information of pixel groups relating to a plurality of polygon graphics having different depths, the graphics in the graphics list (display list) to be displayed are sorted in advance from those closest to the viewer of the screen and the respective graphics are drawn according to the order of the sorting. More specifically, an aspect of the present invention is characterized by the fact that a graphics group in a form of a graphics list containing image information of pixel groups relating to a plurality of polygon graphics is sorted according to the depth information relating to depth of the respective graphics, the respective graphics are selected according to the order of the sorting, whether or not the drawing region of a polygon belonging to the selected graphics is contained in a drawing region of a polygon belonging to a graphic having a higher sorting order is judged according to the image information, a graphics which is judged as contained in the drawing region by the judgment is deleted from the graphics list, and when the drawing for the respective pixels of the graphics judged to be not contained in the drawing region by the judgment is commanded, the drawing for the respective pixels of the graphics is executed based on the image information.\nFurther, an aspect of the present invention is characterized, in that in the course of executing the above graphics drawing method, when a drawing for the respective pixels of the graphics judged as not contained in the drawing region, judgment processing of the graphics list is commanded, the drawing use graphics data of the respective pixels are extracted from the data storage region storing the drawing use graphics data of the respective pixels only when a flag indicating that a pixel is not yet drawn is added in the extracted drawing use graphics data, the drawing for the pixel is executed and a flag indicating that the drawing is completed is added in the drawing use graphics data of the pixel and is stored in the data storage region; and, when a flag indicating that a pixel is already drawn is added in the drawing use graphics data of the respective pixels extracted from the data storage region, the redrawing of the pixel is prohibited.\nAn aspect of the present invention is characterized, by the fact that, when judging specifically the drawing region of a polygon belonging to the graphic list, a circumscribed rectangle for a polygon belonging to the graphics list and an inscribed rectangle for the polygon are respectively provided in correspondence to the drawing region, and whether or not a circumscribed rectangle relating to a selected graphic is hidden by an inscribed rectangle relating to a graphic having a higher order of sorting is judged.\nFurther, an aspect of the present invention is characterized by the fact that a method is used in which a plurality of groups, each containing one image unit consisting of a plurality of polygons, are formed depending on their depth, graphics lists containing image information of the pixel groups corresponding to the polygon graphics are produced for every one of the respective groups, and when judging an overlapping of a graphic belonging to a graphics list, circumscribed rectangles for the graphics of the respective groups and inscribed rectangles for the graphics of the respective groups are respectively provided in correspondence to the drawing region; and, whether or not a circumscribed rectangle relating to the graphic of the selected group is hidden by an inscribed rectangle relating to a graphic of a group having a higher order of sorting is judged.\nFurther, an aspect of the present invention is characterized by the fact that, in order to realize a high speed drawing processing by eliminating a Z comparison in a drawing processing based on a graphics list containing image information of pixel groups corresponding to a plurality of polygon graphics having different depths, the drawing regions for a plurality of graphics belonging to the graphics list are respectively provided, when drawing of the respective graphics is commanded, whether or not the drawing regions relating to the respective graphics belong to a drawing region relating to a graphic having a higher order of drawing is judged and a drawing for a graphic judged as not belonging to the drawing region by the judgment is executed according to the image information.\nAn aspect of the present invention is characterized by the fact that, based on a graphics list containing image information of the pixel groups corresponding to a plurality of polygon graphics having different depths, drawing regions of a plurality of graphics belonging to the graphics list are provided, when drawing for the respective graphics is commanded, whether or not the drawing regions relating to the respective graphics belong to a drawing region relating to a graphic having a higher order of drawing is judged. A drawing of the graphic judged as not belonging to the drawing region by the judgment is executed according to the image information, when a drawing for the respective pixels of the graphic judged as belonging to the drawing region by the judgment is commanded, the drawing use graphics data of the respective pixels are extracted from the data storage region storing the drawing use graphics data of the respective pixels, and only when a flag indicating that a pixel is not yet drawn is added in the extracted drawing use graphics data, will a drawing for the pixel be executed and a flag indicating that the pixel is already drawn is added in the drawing use graphics data of the pixel be added and stored in the data storage region; and, when a flag indicating that the pixel is already drawn is added in the drawing use graphics data of the respective pixels extracted from the data storage region, redrawing of the pixel is prohibited.\nAn aspect of the present invention is characterized by the fact that a plurality of groups, each containing a plurality of polygons as an image unit, are formed, depending on their depths, and a graphics list containing image information of pixel groups corresponding to polygon graphics of the respective groups and circumscribed rectangles for the graphics of the respective groups in the graphics list are respectively provided in correspondence to the drawing regions, when drawing for the graphics of the respective groups is commanded, whether or not the circumscribed rectangles relating to the graphics of the respective groups belong to a circumscribed rectangle relating to a graphic of a group having a higher order of drawing is judged. A drawing for the graphic of the group judged as not belonging is executed according to the image information, and when a drawing for the graphic of the group judged as belonging is commanded, the drawing use graphics data of the respective pixels are extracted from the data storage region storing the drawing use graphics data of the respective pixels of the graphic belonging to the group, and only when a flag indicating that a pixel is not yet drawn is added in the extracted drawing use graphics data, will a drawing for the pixel be executed and a flag indicating that the pixel already drawn is added in the drawing use graphics data of the pixel be added and stored in the data storage region; and, when a flag indicating that the pixel is already drawn is added in the drawing use graphics data of the respective pixels extracted from the data storage region, redrawing of the pixel is prohibited.\nAn aspect of the present invention is characterized by the fact that a method is used in which, in the course of executing a graphics drawing method eliminating Z comparison, in order to judge an overlapping of graphics belonging to a graphics list, circumscribed rectangles for a plurality of polygons are respectively provided in correspondence to the drawing regions, and drawing for the respective polygons is commanded, whether or not the circumscribed rectangles relating to the respective polygons belong to a circumscribed rectangle relating to a polygon having a higher order of drawing.\nAn aspect of the present invention is characterized by a three dimension graphics display method displaying a three dimensional image on a display screen according to the drawing use graphics data obtained through one of the above graphics drawing methods.\nFurther, an aspect of the present invention is characterized by the fact that a graphics drawing device, in which information required for judging overlapping of graphics is reduced, is constituted by a data storage means in which graphics data relating to two dimensional graphics is stored in correspondence to pixel groups, as well as control data indicating whether or not the respective pixels already drawn are stored; a data extracting means which extracts designated graphics data and control data from the data storage means based on the image information of the pixel groups corresponding to a plurality of polygon graphics; a graphics data producing means which, when the control data extracted by the data extracting means indicates that drawing of a pixel is not completed, new graphics data is produced according to the image information, and, when the control data extracted by the data extracting means indicates that drawing of a pixel is completed, the production of new graphics data is eliminated; a drawing means which draws the graphics data produced by the graphics data producing means on the pixel designated by the data storage means; and a control data providing means which provides data representing that the drawing of the pixel is completed for the control data of the pixel which has been drawn by the drawing means among the control data in the data storage means.\nFurther, an aspect of the present invention is characterized by the fact that the above device in which graphics are drawn based on the image information of the graphics list and in which information required for judging an overlapping of graphics is reduced, further comprises a graphics list storage means which stores a graphics list containing the image information of the pixel groups corresponding to a plurality of polygon graphics; a sorting means which sorts graphics groups in the graphics list stored in the graphics list storage means according to depth information relating to the depth of the respective graphics; and a means for transferring the image information relating to the sorted graphics by the sorting means to the data extracting means.\nFurther, an aspect of the present invention is characterized by the fact that a graphics drawing device in which only graphics not overlapping other graphics are extracted among the graphics list, comprises a data storage means in which graphics data relating to two dimensional graphics is stored in correspondence to pixel groups as well as control data indicating whether or not the respective pixels already drawn are stored; a graphics list storage means which stores a graphics list containing image information of the pixel groups corresponding to a plurality of polygon graphics; a sorting means which sorts graphics groups in the graphics list stored in the graphics list storage means according to depth information relating to the depth of the respective graphics; a region judgment means which judges, based on the image information, whether or not the drawing regions of polygons belonging to the sorted graphics by the sorting means are contained in a drawing region of a polygon belonging to a graphic having a higher order of sorting; a deleting means which deletes the graphics judged as contained by the region judgment means from the graphics list; a graphics data producing means which produces graphics data relating to graphics judged as not contained by the region judgment means for every one of the pixels according to the image information; and a drawing means which draws the graphics data produced by the graphics data producing means on the designated pixel in the data storage means.\nAn aspect of the present invention is characterized by the fact that the above graphics drawing device, in which only graphics not overlapping others are extracted from the graphics list, further comprises a rectangle providing means which respectively provides circumscribed rectangles for the polygons belonging to the sorted graphics by the sorting means and inscribed rectangles for the polygons while correlating the data to the drawing regions and a region judgment means which judges the drawing region according to the rectangles provided by the rectangle providing means.\nAn aspect of the present invention is characterized by the fact that the above graphic drawing device in which the graphics list storage means includes a plurality of groups each containing one image unit consisting of a plurality of polygons depending on their depth and stores for every one of the groups graphics list containing the image information of the pixel groups corresponding to the respective polygon graphics, further comprises a rectangle providing means which respectively provides circumscribed rectangles for the graphics of the groups belonging to the sorted graphics by the sorting means and inscribed rectangles for the graphics of the groups while correlating the data to the drawing regions and a region judgment means which judges the drawing region according to the rectangles provided by the rectangle providing means.\nAn aspect of the present invention is characterized by the fact that the graphics drawing device, in which Z comparison processing in the drawing processing is eliminated, comprises a data storage means which stores graphics data relating to two dimensional graphics while correlating the data to the pixel groups, as well as stores control data indicating whether or not the respective pixels are already drawn; graphics list storage means which includes a plurality of groups each containing one image unit consisting of a plurality of polygons depending on their depth and stores for every one of the groups a graphics list containing image information of the pixel groups corresponding to the respective polygon graphics; a sorting means which sorts the graphics groups in the graphics list stored in the graphics list storage means according to depth information relating to the depth of the respective graphics; a rectangle providing means which respectively provides circumscribed rectangles for graphics in a group belonging to the graphics of groups sorted by the sorting means and inscribed rectangles for the graphic in the group while correlating the data to the drawing regions; a region judgment means which judges whether or not the circumscribed rectangle relating to the graphic in the group sorted among the rectangles provided by the rectangle providing means is hidden by an inscribed rectangle belonging to a graphics in a group having a higher order of sorting; a graphics data producing means which produces for every one of the pixels graphics data relating to only the polygons in the graphics judged as not hidden by the region judgment means according to the image information; and a drawing means which draws the graphics data produced by the graphics data producing means on the pixel designated by the data storage means.\nSimilarly, an aspect of the present invention is characterized by the fact that the graphics drawing device comprises a data storage means which stores graphics data relating to two dimensional graphics while correlating the data to the pixel groups and storing control data indicating whether or not the respective pixels are already drawn; a rectangle providing means which, based on graphics list containing the image information of the pixel groups corresponding to a plurality of polygon graphics having different depths, provides circumscribed rectangles for the polygons belonging to the graphics list while correlating the data to the drawing regions; a region judgment means which judges whether or not the respective circumscribed rectangles provided by the rectangle providing means belong to a circumscribed rectangle having a higher order of drawing; a first graphics producing means which produces for every one of the pixels graphics data relating to the polygon judged as not belonging by the region judgment means; a data extracting means which, when judged as belonging by the region judgment means, based on the image information, extracts the graphics data and control data designated by the data storage means; a second graphics data producing means which, when the control data extracted by the data extracting means indicates that the drawing is not completed, produces production of new graphics data according to the image information in the graphics list, and when the control data extracted by the data extracting means indicates that the drawing is already completed, eliminates production of new graphics data; a drawing means which draws the graphics data produced by the respective graphics data producing means on the pixels designated by the data storage means; and a control data providing means which provides data indicating drawing completion for the control data of the pixels drawn by the drawing means among the control data in the data storage means.\nAn aspect of the present invention is characterized by the fact that the respective means in the above respective graphics drawing devices are formed into a one chip LSI and are mounted on a common substrate.\nFurther, an aspect of the present invention is characterized by the fact that a three dimensional graphics display device comprises a display device which displays a three dimensional image on a display screen according to the graphics data of the respective pixels drawn by one of the above drawing devices.\nAn aspect of the present invention is characterized by the fact that, in the course of successively drawing the pixels in the respective graphics only when control data or a flag is stored as information representing that the drawing is not completed in a control data writing region, the drawing is executed according to the image information, and when control data or a flag is stored as information representing that the drawing is completed in the control data writing region, the redrawing is prohibited, so that it is enough only to judge the information of the control data writing region for judging an overlapping of the graphics, with the result that the amount of information needed for judging the overlapping of the graphics can be limited and the capacity of the data storage region storing the graphics data can be reduced.\nAn aspect of the present invention is characterized by the fact that, in the course of drawing respective graphics based on the graphics list, when a drawing region of a polygon belonging to the graphics list is contained in a drawing region of a polygon belonging to a graphic having a higher order of drawing, the former graphic is deleted from the graphics list and only a drawing for the respective pixel of the graphics judged as not contained is performed, so that only graphics not overlapping others are extracted and a high speed drawing processing is achieved.\nFurther, an aspect of the present invention is characterized by the fact that, in the course of drawing the respective graphics, it is judged whether or not drawing regions of the respective graphics belong to a drawing region relating to a graphic having a higher order of drawing and drawing for the graphics judged as not belonging is executed according to the image information, such that Z comparison in the drawing processing is eliminated and a high speed drawing processing is achieved.\nAn aspect of the present invention is characterized by the fact that in order to realize a high speed image processing system through parallel processing and use of a high speed memory access in the system, the processing performance of a CPU is not reduced even when a processor other than the CPU accesses the memory. Namely, in order to permit a plurality of processors to simultaneously access a memory, a bus use bit range for the respective processors is controlled.\nFurther, an aspect of the present invention is characterized by the fact that, in order for every module to control a memory independently, at least one control line can be controlled separately.\nFurther, an aspect of the present invention is characterized by the fact that, in order to permit respective memory modules to access different row addresses, a control which permits a row address activating command to be selectively set in a memory is performed.\nFurther, an aspect of the present invention is characterized by the fact that high quality image processing arithmetic in the image input and output processing is performed through a spacial filtering while making use of multi-valued information, coding and decoding processing arithmetic is performed through table retrieval based on variation point detection and a bit pattern, modern arithmetic in the communication processing is performed through digital filtering while making use of sum of product arithmetic. A one chip processor performing the above processings and constituted by a high speed RISC and an image processing core and a high speed memory are connected via a dedicated bus, and the processings and memory are integrated.\nFurther, an aspect of the present invention is characterized by the fact that, memories for the respective processings are mapped on a common memory address.\nFurther, an aspect of the present invention is characterized by the fact that, in order to execute the respective processings in a multi-task manner, a register is provided which dynamically allocates bus occupying time for the respective processings depending on their predetermined processing loads so as not to drop below a minimum bus throughput for guaranteeing a real time operation for the respective processings. The register can be easily modified by software, and the set values can be modified, such as by the processing sequences and the operating modes.\nFurther, an aspect of the present invention is characterized by the fact that, in order to enhance bus efficiency, a means is provided which predicts a change in the row address and issues in advance an address control command.\nFurther, an aspect of the present invention is characterized by the fact that, in order to reset effective data in an internal arithmetic register at the time of task switching, a means is provided which returns a read address and executes restarting of a processing.\nAn aspect of the present invention is characterized by the fact that, through control of the bus bit range by the processor, if no conflicts between bit ranges are caused, simultaneous access can be achieved.\nAn aspect of the present invention is characterized by the fact that, since the memory bus can be adjusted by the bit range, the CPU and the image processing use a processor which can access the memory in parallel, depending on the operation mode, and the processing can be continued with no overhead for mishits and high speed processing is achieved."}
{"text": "Side-by-side vehicles are known. The present disclosure relates to vehicles, including utility vehicles. The present disclosure relates to air handling systems for vehicles. The present disclosure relates to suspension systems for vehicles.\nIn exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle, comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including seating and operator controls; a rear drive unit supported by the frame and positioned rearward of the operator area; and a first rear suspension system moveably coupling the first ground engaging member to the frame. The rear drive being operatively coupled to the power source and operatively coupled to at least a first ground engaging member positioned rearward of the operator area through a drive shaft to transfer power received from the power source to the first ground engaging member. The first rear suspension system including a control arm moveably coupled to the frame at a location rearward of the drive shaft connecting the rear drive unit and the first ground engaging member and between a plane passing through a first laterally extending end of the rear drive unit and a centerline longitudinal plane of the vehicle. The plane being parallel to the centerline longitudinal plane of the vehicle. In one example, the first rear suspension includes a second control arm moveably coupled to the frame at a second location rearward of the drive shaft connecting the rear drive unit and the first ground engaging member and between the plane passing through the first laterally extending end of the rear drive unit and the centerline longitudinal plane of the vehicle. In a variation thereof, the control arm and the second control arm are coupled to a bearing carrier which is coupled to the first ground engaging member, the bearing carrier including an opening through which the drive shaft is operatively coupled to the first ground engaging member. In a further variation thereof, the first rear suspension includes a radius arm coupled to the frame at a third location forward of the rear drive unit and coupled to the bearing carrier. The third location may be forward of the power source. The third location may be under the seating of the operator area. In another variation thereof, the first rear suspension includes a dampening member having a first end coupled to the radius arm at a fourth location and a second end coupled to the frame at a fifth location, the fifth location being forward of the fourth location and closer to the longitudinal centerline plane than the fourth location. In another example, the vehicle further comprises a front drive unit supported by the frame and positioned forward of the operator area. The front drive being operatively coupled to the power source and operatively coupled to at least a second ground engaging member positioned forward of the operator area to transfer power received from the power source to the second ground engaging member.\nIn another exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including seating and operator controls; a rear drive unit supported by the frame and positioned rearward of the operator area, the rear drive being operatively coupled to the power source and operatively coupled to at least a first ground engaging member positioned rearward of the operator area through a drive shaft to transfer power received from the power source to the first ground engaging member; and a first rear suspension system moveably coupling the first ground engaging member to the frame. The first rear suspension system including a control arm coupled to the frame at a rear face of the frame. The control arm being unobscured by the frame from a viewing direction which is perpendicular to a centerline longitudinal plane of the vehicle and rearward of the vehicle. In one example, the rear face of the vehicle is a flat surface. In another example, the first rear suspension includes a second control arm moveably coupled to the rear face of the frame at a second location the second control arm being unobscured by the frame from the viewing direction which is perpendicular to a centerline longitudinal plane of the vehicle and rearward of the vehicle. In a variation thereof, the control arm and the second control arm are coupled to a bearing carrier which is coupled to the first ground engaging member. The bearing carrier including an opening through which a drive shaft is operatively coupled to the first ground engaging member. The first rear suspension may include a radius arm coupled to the frame at a third location forward of the power source and coupled to the bearing carrier. The third location may be under the seating of the operator area. The first rear suspension includes a dampening member having a first end coupled to the radius arm at a fourth location and a second end coupled to the frame at a fifth location. The fifth location being forward of the fourth location and closer to the longitudinal centerline plane than the fourth location. In another example, the vehicle further comprises a front drive unit supported by the frame and positioned forward of the operator area, the front drive being operatively coupled to the power source and operatively coupled to at least a second ground engaging member positioned forward of the operator area to transfer power received from the power source to the second ground engaging member.\nIn another exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including seating and operator controls; a rear drive unit supported by the frame and positioned rearward of the operator area. The rear drive being operatively coupled to the power source and operatively coupled to a first ground engaging member positioned rearward of the operator area through a first drive shaft to transfer power received from the power source to the first ground engaging member and to a second ground engaging member positioned rearward of the operator area through a second drive shaft to transfer power received from the power source to the second ground engaging member. The first ground engaging member being positioned on a first side of a vertical centerline longitudinal plane of the vehicle and the second ground engaging member being positioned on a second side of the vertical centerline longitudinal plane of the vehicle. The vehicle further comprising a first rear suspension system moveably coupling the first ground engaging member to the frame. The first rear suspension system including a first radius arm coupled to the first ground engaging member and coupled to the frame at a first location forward of the first ground engaging member, a first control arm coupled to the first ground engaging member and to the frame; and a first dampening member coupled to the first radius arm and to the frame. The vehicle further comprising a second rear suspension system moveably coupling the second ground engaging member to the frame. The second rear suspension system including a second radius arm coupled to the second ground engaging member and coupled to the frame at a second location forward of the second ground engaging member, a second control arm coupled to the second ground engaging member and to the frame; and a second dampening member coupled to the second radius arm and to the frame. The vehicle further comprising a sway bar coupling the first rear suspension to the second rear suspension, the sway bar being coupled to the frame at a location forward of the power source. In one example, the location at which the sway bar is coupled to the frame is rearward of the first location and the second location. In another example, the sway bar is coupled to the first radius arm through a first link and the sway bar is coupled to the second radius arm through a second link. In a variation thereof, the location at which the sway bar is coupled to the frame is rearward of the first location and the second location. In another example, the first control arm is coupled to the frame at a fourth location, the power source being positioned between the first location and the fourth location. In a variation thereof, the location at which the sway bar is coupled to the frame is rearward of the first location and the second location.\nIn a further exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including side-by-side seating and operator controls; a rear drive unit supported by the frame and positioned rearward of the operator area; and a first rear suspension system moveably coupling the first ground engaging member to the frame. The rear drive being operatively coupled to the power source and operatively coupled to a first ground engaging member positioned rearward of the operator area through a first drive shaft to transfer power received from the power source to the first ground engaging member and to a second ground engaging member positioned rearward of the operator area through a second drive shaft to transfer power received from the power source to the second ground engaging member, the first ground engaging member being positioned on a first side of a vertical centerline longitudinal plane of the vehicle and the second ground engaging member being positioned on a second side of the vertical centerline longitudinal plane of the vehicle. The first rear suspension system including a first moveable arm. The vehicle further comprising a second rear suspension system moveably coupling the second ground engaging member to the frame, the second rear suspension system including a second moveable arm; and a sway bar coupling the first rear suspension to the second rear suspension, the sway bar being coupled to the frame at a location forward of the power source. In one example, the sway bar is coupled to the first moveable arm and the second moveable arm. In a variation thereof, the vehicle further comprises a front drive operatively coupled to the power source and operatively coupled to a third ground engaging member positioned forward of the operator area and a fourth ground engaging member positioned forward of the operator area, the third ground engaging member being positioned on the first side of the vertical centerline longitudinal plane of the vehicle and the fourth ground engaging member being positioned on the second side of the vertical centerline longitudinal plane of the vehicle.\nIn a further exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including seating and operator controls; a rear drive unit supported by the frame and positioned rearward of the operator area, the rear drive being operatively coupled to the power source and operatively coupled to a first ground engaging member positioned rearward of the operator area through a first drive shaft to transfer power received from the power source to the first ground engaging member; and a first rear suspension system moveably coupling the first ground engaging member to the frame. The first rear suspension system including a first radius arm coupled to the first ground engaging member and coupled to the frame at a first location forward of the first ground engaging member, a first control arm coupled to the first ground engaging member and to the frame; and a first dampening member coupled to the first radius arm and to the frame. The vehicle further comprising a bearing carrier coupled to the first ground engaging member, the first radius arm, and the first control arm. The bearing carrier including an opening through which the drive shaft is operatively coupled to the first ground engaging member. The first radius arm may be uncoupled from the frame at the first location and uncoupled from the bearing carrier without uncoupling the drive shaft from the first ground engaging member. In one example, the first radius arm includes an open end which receives the drive shaft.\nIn still another exemplary embodiment of the present disclosure, a method of removing a radius arm of a rear suspension of a vehicle is provided. The method comprising the steps of (a) uncoupling a first portion of the radius arm of the rear suspension from a frame of the vehicle; (b) uncoupling a second portion of the radius arm of the rear suspension from a bearing carrier which is coupled to a wheel of the vehicle, the bearing carrier having an opening through which a drive shaft is operatively coupled to the ground engaging member; and (c) removing the radius arm from the vehicle, wherein the drive shaft remains coupled to the wheel throughout steps (a) through (c). In one example, the method further comprises the step of uncoupling a third portion of the radius arm of the rear suspension from the frame of the vehicle, the third portion being spaced apart from the first portion and the second portion. In a variation thereof, the step of uncoupling the third portion of the radius arm of the rear suspension from the frame of the vehicle includes the step of uncoupling the third portion of the radius arm of the rear suspension from a dampening member which is coupled to the frame.\nIn yet a further exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; and a unit. The unit including a power source supported by the frame through less than three connections, a transmission supported by the frame through less than three connections, and a spacer coupled to the power source housing and coupled to the transmission housing to position the transmission relative to the power source. The power source having a power source housing and a power source output drive member. The transmission having a transmission housing and a transmission input drive member and a transmission output drive member. The transmission input drive member being operatively coupled to the power source output drive member and the transmission output drive member being operatively coupled to at least one of the plurality of the plurality of ground engaging members to propel the vehicle. The power source output drive member and the transmission input drive member being completely outside of the spacer. The unit being supported by the frame through at least a first connection, a second connection, and a third connection. In one example, the spacer is coupled to the power source through a first fastener and a second fastener and the spacer is coupled to the transmission through a third fastener and a fourth fastener. In a variation thereof, the first fastener is parallel to the second fastener and the third fastener is parallel to the fourth fastener. In another example, the spacer is received by a portion of the power source housing and is coupled to the power source housing through a first connection and a second connection and wherein the spacer is received by a portion of the transmission housing and is coupled to the transmission through a third connection and a fourth connection.\nIn a further exemplary embodiment of the present disclosure, a method of supporting a power source and a transmission on a frame of a vehicle is provided. The method comprising the steps of coupling the power source to the transmission with a spacer, a power source output drive member of the power source and an input drive member of the transmission being completely outside of the spacer; and supporting the power source, the transmission, and the spacer on the frame through at least three connections, less than three connections supporting the power source and less than three connection supporting the transmission.\nIn still a further exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including seating and operator controls; a rear drive unit supported by the frame and positioned rearward of the operator area, the rear drive being operatively coupled to the power source and operatively coupled to at least a first ground engaging member positioned rearward of the operator area through a drive shaft to transfer power received from the power source to the first ground engaging member; and a first rear suspension system moveably coupling the first ground engaging member to the frame. The first rear suspension system including a radius arm moveably coupled to the frame at a first location forward of the drive shaft; a control arm moveably coupled to the frame at second location rearward of the drive shaft; and a dampening member moveably coupled to the frame at a third location and moveably coupled to the radius arm at a fourth location. The third location being forward of the fourth location and closer to the longitudinal centerline plane than the fourth location. In one example, the vehicle further comprises a bearing carrier coupled to the control arm and the radius arm, the first suspension being generally rotatable about a first suspension axis generally passing through the first location and the second location. A top view projection of the first suspension axis intersecting a longitudinal centerline of the vehicle. In a variation thereof, a top view projection of a centerline of the dampening member is angled relative to the top view projection of the first suspension axis. The top view projection of the centerline of the dampening member may be generally normal to the top view projection of the first suspension axis. The top view projection of the centerline of the dampening member may be angled up to about 30 degrees from a normal to the top view projection of the first suspension axis. The top view projection of the centerline of the dampening member may be angled up to about 20 degrees from a normal to the top view projection of the first suspension axis. The top view projection of the centerline of the dampening member may be angled up to about 10 degrees from a normal to the top view projection of the first suspension axis.\nIn yet another exemplary embodiment, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an operator area supported by the frame, the operator area including seating and operator controls; a cargo carrying portion supported by the frame and located rearward of the operator area; and a first air intake system operatively coupled to the power source to communicate ambient air to the power source. The first air intake system receiving ambient air through an inlet in an exterior body panel of the cargo carrying portion. In one example, the cargo carrying portion is a cargo bed. In a variation thereof, the cargo carrying portion includes a floor and a plurality of walls. The floor including a removable cover which permits access to a portion of the first air intake system. The portion of the first air intake system may be an airbox including a filter. In another example, the first air intake system includes a resonator box located between an exterior surface of the exterior body panel of the cargo carrying portion and a wall of a cargo carrying region of the cargo carrying portion. In yet another example, the vehicle further comprises a cover coupled to the exterior body panel to cover the inlet in the exterior body panel. The cover permitting the ambient air to pass by the cover. In a variation thereof, the vehicle further comprises a filter housing positioned behind the cover and a filter removably positioned within the filter housing. The ambient air passing through the filter. The filter may be located between an exterior surface of the exterior body panel of the cargo carrying portion and a wall of a cargo carrying region of the cargo carrying portion. The inlet in the exterior body panel of the cargo carrying portion may be positioned on a first side of a vertical centerline plane of the vehicle and the first air intake system transports the ambient air to a second side of the vertical centerline plane as the ambient air travels through a fluid conduit of the first air intake system. The vehicle may further comprise a CVT unit supported by the frame and operatively coupled between the power source and the at least one of the plurality of ground engaging members; and a second air intake system operatively coupled to the CVT unit to communicate ambient air to an interior of the CVT unit. The second air intake system including a second air inlet through which ambient air enters the second air intake system. The second air inlet being completely positioned to the second side of the vertical centerline plane of the vehicle. The second air intake system transporting the ambient air received through the second air inlet to the first side of the vertical centerline plane of the vehicle as the ambient air travels towards the CVT unit.\nIn still another exemplary embodiment of the present disclosure, a method of providing air to a power source of a vehicle is provided. The method comprising the steps of: providing an air inlet of an air intake system in a cargo carrying portion of the vehicle, the air inlet being rearward of an operator area of the vehicle; receiving a first amount of ambient air through the air inlet; and communicating the first amount of ambient air to the power source of the vehicle. In one example, the air inlet is provided in an exterior body panel of the cargo carrying portion of the vehicle. In another example, the cargo carrying portion includes a cargo bed. In a further example, the step of communicating the first amount of ambient air to the power source of the vehicle includes the step of: passing the ambient air through a first filter located between an exterior of the cargo carrying portion and a wall of a cargo bed of the cargo carrying portion. In still another example, the step of communicating the first amount of ambient air to the power source of the vehicle includes the step of passing the ambient air through a resonator box located between an exterior of the cargo carrying portion and a wall of a cargo bed of the cargo carrying portion. In yet still another example, the step of communicating the first amount of ambient air to the power source of the vehicle includes the step of passing the ambient air through an airbox. The airbox including a base portion, a cover portion, and a filter positioned in an interior of the airbox. The ambient air passing through the filter, the airbox being located below a floor of the cargo carrying portion. In still another example, the method further comprises the steps of manipulating a portion of the floor of the cargo carrying portion to access the airbox from above the floor of the cargo carrying portion; and moving the cover portion of the airbox relative to the base portion of the airbox to access the filter of the airbox, the cover portion being rotatably coupled to the base portion.\nIn yet still another exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; a CVT unit supported by the frame and operatively coupled between the power source and the at least one of the plurality of ground engaging members; an operator area supported by the frame, the operator area including seating and operator controls; a cargo carrying portion supported by the frame and located rearward of the operator area; and an air intake system operatively coupled to the CVT unit to communicate ambient air to an interior of the CVT unit, the air intake system receiving ambient air through an inlet in an exterior body panel of the cargo carrying portion. In one example, the cargo carrying portion is a cargo bed. In another example the vehicle further comprises a cover coupled to the exterior body panel to cover the inlet in the exterior body panel. The cover permitting the ambient air to pass by the cover. In a further example, the vehicle further comprises a filter housing positioned behind the cover and a filter removably positioned within the filter housing. The ambient air passing through the filter. In a variation thereof, the filter is located between an exterior surface of the exterior body panel of the cargo carrying portion and a wall of a cargo carrying region of the cargo carrying portion. The inlet in the exterior body panel of the cargo carrying portion may positioned on a first side of a vertical centerline plane of the vehicle and the air intake system transports the ambient air to a second side of the vertical centerline plane as the ambient air travels through a fluid conduit of the air intake system. In another example, the CVT unit includes a drive member operatively coupled to the power source; a driven member operatively coupled to the at least one ground engaging member; a drive belt operatively coupling the driven member to the drive member; and a CVT housing having an interior containing the drive member, the driven member, and the drive belt. The CVT housing including a plurality of air inlets to the interior of the CVT housing, a first air inlet being positioned proximate the drive member and a second air inlet being positioned proximate the driven member. Both the first air inlet and the second air inlet being in fluid communication with the air intake system to receive ambient air from the air intake system. In a variation thereof, the CVT housing includes an air outlet through which air exits the interior of the CVT housing, the air outlet being in fluid communication with a fluid conduit which directs the air at a portion of the power source.\nIn still another exemplary embodiment of the present disclosure, a method of providing air to a CVT unit of a vehicle is provided. The method comprising the steps of: providing an air inlet of an air intake system in a cargo carrying portion of the vehicle, the air inlet being rearward of an operator area of the vehicle; receiving a first amount of ambient air through the air inlet; and communicating the first amount of ambient air to an interior of the CVT unit of the vehicle. In one example, the air inlet is provided in an exterior body panel of the cargo carrying portion of the vehicle. In another example, the cargo carrying portion includes a cargo bed. In still another example, the step of communicating the first amount of ambient air to the CVT unit of the vehicle includes the step of: passing the ambient air through a first filter located between an exterior of the cargo carrying portion and a wall of a cargo bed of the cargo carrying portion.\nIn yet another exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; an operator area supported by the frame, the operator area including seating and operator controls; a cargo carrying portion supported by the frame and located rearward of the operator area; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; a CVT unit supported by the frame at a first position; a first air intake system operatively coupled to the power source to communicate ambient air to the power source, the first air intake system including a first air inlet through which ambient air enters the first air intake system, the first air inlet being completely positioned to a first side of the CVT unit; and a second air intake system operatively coupled to the CVT unit to communicate ambient air to an interior of the CVT unit, the second air intake system including a second air inlet through which ambient air enters the second air intake system, the second air inlet being completely positioned to a second side of the CVT unit. The CVT unit being operatively coupled between the power source and the at least one of the plurality of ground engaging members. In one example, the first air inlet and the second air inlet are both rearward of a front plane of the seating of the operator area. In a variation thereof, the first air inlet and the second air inlet are both rearward of the operator area. In another variation thereof the vehicle further comprises a roll cage. The first air inlet and the second air inlet are both rearward of the roll cage. In another example, the first air inlet and the second air inlet are positioned above the plurality of ground engaging members. In a variation thereof the vehicle further comprises a rear drive unit supported by the frame rearward of the front plane of the seating and operatively coupled to at least one ground engaging member which is rearward of the front plane of the seating, the power source being operatively coupled to the rear drive unit; and a first suspension system coupling the at least one ground engaging member which is rearward of the front plane of the seating to the frame, the first air inlet and the second air inlet being positioned completely above the first suspension system.\nIn still another exemplary embodiment of the present disclosure, a vehicle is provided. The vehicle comprising a frame; a plurality of ground engaging members supporting the frame; a power source supported by the frame and operatively coupled to at least one of the plurality of ground engaging members to propel the vehicle; an air intake system operatively coupled to the CVT unit to communicate ambient air to an interior of the CVT unit; and a CVT unit supported by the frame. The CVT unit being operatively coupled between the power source and the at least one of the plurality of ground engaging members. The CVT unit comprising a drive member operatively coupled to the power source; a driven member operatively coupled to the at least one ground engaging member; a drive belt operatively coupling the driven member to the drive member; and a CVT housing having an interior containing the drive member, the driven member, and the drive belt. The CVT housing including a plurality of air inlets to the interior of the CVT housing. A first air inlet being positioned proximate the drive member and a second air inlet being positioned proximate the driven member. Both the first air inlet and the second air inlet being in fluid communication with the air intake system to receive ambient air from the air intake system. In one example, the CVT unit includes a diverter which receives the ambient air from the air intake system and directs a first portion of the ambient air to the first air inlet and a second portion of the ambient air to the second air inlet. In a variation thereof, the CVT housing includes a base portion and a cover. The cover being removably coupled to the base portion. The diverter being associated with the base portion. The diverter may be part of the base portion of the CVT housing. In another example, the first portion of the ambient air moves from the first air inlet towards an air outlet of the CVT housing in a generally counterclockwise movement and the second portion of the ambient air moves from the second air inlet towards the air outlet of the CVT housing in a generally counterclockwise movement.\nThe above mentioned and other features of the invention, and the manner of attaining them, will become more apparent and the invention itself will be better understood by reference to the following description of embodiments of the invention taken in conjunction with the accompanying drawings. These above mentioned and other features of the invention may be used in any combination or permutation.\nCorresponding reference characters indicate corresponding parts throughout the several views. Unless stated otherwise the drawings are proportional."}
{"text": "Fiber optic cables are being proposed for use by the military. These applications include long length communications systems, computer interconnects, data links for remote operations, and local distribution systems. Many of these applications will require the cables to function properly in nuclear radiation environments in which high levels of thermal radiation can be encountered. A considerable body of literature exists showing the deleterious effects of ionizing radiation on optical waveguides. However, little attention has been given to protecting the optical waveguides from the effects of thermal radiation pulses which accompany the ionizing radiation.\nIt is therefore an object of this invention to provide means to mitigate damage to optical fiber cables which occurs upon exposure of the cables to thermal radiation.\nIt is an object of this invention to provide such protective means for the cables which is relatively thin and lightweight, and which adds little to the cost of the cables."}
{"text": "As graphical user interfaces and display systems for computers continue to offer increased functionality, they provide the user with greater flexibility in manipulating the information that is displayed on a monitor, or similar type of display device. In some instances, however, such manipulation could produce unexpected results. For example, some display systems offer the user the ability to change the resolution of the display device. If the user decreases the resolution of the device, for example to make icons or text appear larger, the objects on the display will move relative to a reference point, as well as change size, in accordance with the change in resolution. The user may prefer that certain objects remain in their original position, however. For instance, it may be desirable to keep tool palettes in a corner of the display. Furthermore, some of the objects at the periphery of the original display may be partially removed from view, even though the user may desire to have such objects remain readily accessible. Such objects might include utility windows that the user may want to frequently access during the operation of the computer.\nSimilar situations can occur when objects are moved on the desktop of the user interface, or when display configurations change. For instance, if the user switches the display from a full-sized monitor to an LCD screen on a notebook computer, the different operating parameters of the two display devices may affect the size of the desktop. As a result, the positions and sizes of the objects on the display can change, and some of them may become unviewable because they are positioned outside of the display area. It is desirable, therefore, to provide a mechanism which is capable of maintaining windows and other similar objects in an accessible condition that conforms to the user's preferences, despite changes in the display environment."}
{"text": "In recent years, research and development of light-emitting elements using electroluminescence (EL) have been extensively conducted. In the basic structure of such a light-emitting element, a light-emitting substance is interposed between a pair of electrodes. By applying a voltage to this element, light emission can be obtained from the light-emitting substance.\nSince this type of light-emitting element is a self-luminous type, it has advantages over a liquid crystal display in that visibility of a pixel is high and that no backlight is needed. Therefore, light-emitting elements are thought to be suitable as flat panel display elements. Further, such a light-emitting element also has advantages in that the element can be fabricated to be thin and lightweight and that response speed is very high.\nFurther, since this type of a light-emitting element can be formed to have a film shape, surface light emission can be easily obtained. This feature is difficult to realize with point light sources typified by a filament lamp and an LED or with linear light sources typified by a fluorescent light. Therefore, such light-emitting elements also have a high utility value as surface light source that can be applied to lighting apparatuses or the like.\nLight-emitting elements using electroluminescence are broadly classified according to whether they use an organic compound or an inorganic compound as a light-emitting substance.\nWhen an organic compound is used as a light-emitting substance, by application of a voltage to a light-emitting element, electrons and holes are injected into a layer containing the light-emitting organic compound from a pair of electrodes, whereby a current flows. Then, the carriers (i.e., electrons and holes) recombine to place the light-emitting organic compound into an excited state. The light-emitting organic compound returns to a ground state from the excited state, thereby emitting light. Such a light-emitting element using a light-emitting organic compound is generally referred to as an organic EL element.\nNote that an excited state of an organic compound can be of two types: a singlet excited state and a triplet excited state, and luminescence from the singlet excited state is referred to as fluorescence, and luminescence from the triplet excited state is referred to as phosphorescence.\nProposal of a heterostructure in which layers of different organic compounds are stacked has brought about significant development of such light-emitting elements (see Non-Patent Document 1). That is because adopting a heterostructure increases carrier recombination efficiency and then improves emission efficiency. Then, a double heterostructure in which a light-emitting layer is formed between a hole-transport layer and an electron-transport layer has been proposed (see Non-Patent Document 2). An element structure based on the double heterostructure is now mainstream.\nFurther, by increasing the number of stacked layers, an attempt to improve lifetime has been made lately. For example, in Patent Document 1, an increase in lifetime is attempted by providing a layer to decelerate transport of electrons in addition to an electron-transport layer.    [Patent Document 1] Japanese Published Patent Application No. 2006-66890.    [Non-Patent Document 1] C. W. Tang and S. A. Vanslyke, Applied Physics Letters, Vol. 51, No. 12, 913-915, 1987.    [Non-Patent Document 2] C. Adachi, et al., Japanese Journal of Applied Physics, Vol. 27, L269-L271, 1988."}
{"text": "The present invention relates to the uses of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins as therapeutic or prophylactic agents. The substances are active as wound healing, anti-bacterial and/or anti-inflammatory agents.\nEnamel matrix proteins such as those present in enamel matrix are most well-known as precursors to enamel. Enamel proteins and enamel matrix derivatives have previously been described in the patent literature to induce hard tissue formation (i.e. enamel formation, U.S. Pat. No. 4,672,032 (Slavkin)) or binding between hard tissues (EP-B-0 330 967 and EP-B-0 263 086). Thus, the prior art is solely centred on regeneration of hard tissues, while the present application deals with beneficial effects on soft tissue wound healing and anti-bacterial and anti-inflammatory effects which are unexpected findings.\nThe present invention is based on the finding that enamel matrix, enamel matrix derivatives and/or enamel matrix proteins (the term xe2x80x9can active enamel substancexe2x80x9d is in the following also used for an enamel matrix, an enamel matrix derivative or an enamel matrix protein) are beneficial agents for the enhancement or improvement of the healing of wounds in soft tissues (i.e. non-mineralised tissues) such as collagen or epithelium containing tissues, including skin and mucosa, muscle, blood and lymph vessels, nerve tissues, glands, tendons, eyes and cartilage. As demonstrated in the experimental section herein, the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins exert especially useful effects in the healing or prophylaxis of soft tissue wounds.\nAccordingly, the invention relates to the use of a preparation of an active enamel substance for the preparation of a pharmaceutical or cosmetic composition i) for healing of a wound, ii) for improving healing of a wound, and/or iii) for soft tissue regeneration and/or repair.\nIn another aspect, the invention relates to a method of improving the healing of a wound or of promoting soft tissue regeneration and/or repair, the method comprising administering, to an individual in need thereof, a therapeutically or prophylactically effective amount of an active enamel substance.\nFurthermore, enamel matrix, enamel matrix derivatives and enamel matrix proteins have been found to have anti-bacterial and/or anti-inflammatory properties that can be used for treatment of both soft and hard (i.e. mineralised) tissue conditions.\nIn other aspects the invention relate to the use of a preparation of an active enamel substance for the preparation of a pharmaceutical composition for the prevention and/or treatment of an infection or an inflammatory condition.\nWound Healing\nWounds and/or ulcers are normally found protruding from the skin or on a mucosal surface or as a result of an infarction in an organ (xe2x80x9cstrokexe2x80x9d). A wound may be a result of a soft tissue defect or a lesion or of an underlying condition. Regeneration of experimentally provoked periodontal wounds has previously been described by the inventors and is not intended to be within the scope of the present invention. In the present context the term xe2x80x9cskinxe2x80x9d relates to the outermost surface of the body of an animal including a human and embraces intact or almost intact skin as well as an injured skin surface. The term xe2x80x9cmucosaxe2x80x9d relates to undamaged or damaged mucosa of an animal such as a human and may be the oral, buccal, aural, nasal, lung, eye, gastrointestinal, vaginal, or rectal mucosa.\nIn the present context the term xe2x80x9cwoundxe2x80x9d denotes a bodily injury with disruption of the normal integrity of tissue structures. The term is also intended to encompass the terms xe2x80x9csorexe2x80x9d, xe2x80x9clesionxe2x80x9d, xe2x80x9cnecrosisxe2x80x9d and xe2x80x9culcerxe2x80x9d. Normally, the term xe2x80x9csorexe2x80x9d is a popular term for almost any lesion of the skin or mucous membranes and the term xe2x80x9culcerxe2x80x9d is a local defect, or excavation, of the surface of an organ or tissue, which is produced by the sloughing of necrotic tissue. Lesion generally relates to any tissue defect. Necrosis is related to dead tissue resulting from infection, injury, inflammation or infarctions.\nThe term xe2x80x9cwoundxe2x80x9d used in the present context denotes any wound (see below for a classification of wounds) and at any particular stage in the healing process including the stage before any healing has initiated or even before a specific wound like a surgical incision is made (prophylactic treatment).\nExamples of wounds which can be prevented and/or treated in accordance with the present invention are, e.g., aseptic wounds, contused wounds, incised wounds, lacerated wounds, non-penetrating wounds (i.e. wounds in which there is no disruption of the skin but there is injury to underlying structures), open wounds, penetrating wound, perforating wounds, puncture wounds, septic wounds, subcutaneous wounds, etc. Examples of sores are bed sores, canker sores, chrome sores, cold sores, pressure sores etc. Examples of ulcers are, e.g., peptic ulcer, duodenal ulcer, gastric ulcer, gouty ulcer, diabetic ulcer, hypertensive ischemic ulcer, stasis ulcer, ulcus cruris (venous ulcer), sublingual ulcer, submocous clear, symptomatic ulcer, trophic ulcer, tropical ulcer, veneral ulcer, e.g. caused by gonorrhoea (including urethritis, endocervicitis and proctitis). Conditions related to wounds or sores which may be successfully treated according to the invention are burns, anthrax, tetanus, gas gangrene, scalatina, erysipelas, sycosis barbae, folliculitis, impetigo contagiosa, or impetigo bullosa, etc. There is often a certain overlap between the use of the terms xe2x80x9cwoundxe2x80x9d and xe2x80x9culcerxe2x80x9d and xe2x80x9cwoundxe2x80x9d and xe2x80x9csorexe2x80x9d and, furthermore, the terms are often used at random. Therefore as mentioned above, in the present context the term xe2x80x9cwoundxe2x80x9d encompasses the term xe2x80x9culcerxe2x80x9d, xe2x80x9clesionxe2x80x9d, xe2x80x9csorexe2x80x9d and xe2x80x9cinfarctionxe2x80x9d, and the terms are indiscriminately used unless otherwise indicated.\nThe kinds of wounds to be treated according to the invention include also i) general wounds such as, e.g., surgical, traumatic, infectious, ischemic, thermal, chemical and bullous wounds; ii) wounds specific for the oral cavity such as, e.g., post-extraction wounds, endodontic wounds especially in connection with treatment of cysts and abscesses, ulcers and lesions of bacterial, viral or autoimmunological origin, mechanical, chemical, thermal, infectious and lichenoid wounds; herpes ulcers, stomatitis aphthosa, acute necrotising ulcerative gingivitis and burning mouth syndrome are specific examples; and iii) wounds on the skin such as, e.g., neoplasm, burns (e.g. chemical, thermal), lesions (bacterial, viral, autoimmunological), bites and surgical incisions. Another way of classifying wounds is as i) small tissue loss due to surgical incisions, minor abrasions and minor bites, or as ii) significant loss. The latter group includes ischemic ulcers, pressure sores, fistulae, lacerations, severe bites, thermal burns and donor site wounds (in soft and hard tissues) and infarctions.\nThe healing effect of an active enamel substance has been found to be of interest in connection with wounds which are present in the oral cavity. Such wounds may be bodily injuries or trauma associated with oral surgery including periodontal surgery, tooth extraction(s), endodontic treatment, insertion of tooth implants, application and use of tooth prothesis, and the like. In the experimental section herein the beneficial effect of an active enamel substance on such wounds has been demonstrated. Furthermore, a soft tissue effect has been observed.\nIn the oral cavity healing of wounds like aphthous wounds, traumatic wounds or herpes associated wounds is also improved after application of an active enamel substance. The traumatic wounds and the herpes associated wounds can of course also be situated on other parts of the body than in the oral cavity.\nIn other aspects of the invention, the wound to be prevented and/or treated is selected from the group consisting of aseptic wounds, infarctions, contused wounds, incised wounds, lacerated wounds, non-penetrating wounds, open wounds, penetrating wounds, perforating wounds, puncture wounds, septic wounds and subcutaneous wounds.\nOther wounds which are of importance in connection with the present invention are wounds like ischemic ulcers, pressure sores, fistulae, severe bites, thermal burns and donor site wounds.\nIschemic ulcers and pressure sores are wounds which normally only heal very slowly and especially in such cases an improved and more rapid healing is of course of great importance for the patient. Furthermore, the costs involved in the treatment of patients suffering from such wounds are markedly reduced when the healing is improved and takes place more rapidly.\nDonor site wounds are wounds which e.g. occur in combination with removal of hard tissue from one part of the body to another part of the body e.g. in connection with transplantation. The wounds resulting from such operations are very painful and an improved healing is therefore most valuable.\nThe term xe2x80x9cskinxe2x80x9d is used in a very broad sense embracing the epidermal layer of the skin andxe2x80x94in those cases where the skin surface is more or less injuredxe2x80x94also the dermal layer of the skin. Apart from the stratum corneum, the epidermal layer of the skin is the outer (epithelial) layer and the deeper connective tissue layer of the skin is called the dermis.\nSince the skin is the most exposed part of the body, it is particularly susceptible to various kinds of injuries such as, e.g., ruptures, cuts, abrasions, burns and frostbites or injuries arising from the various diseases. Furthermore, much skin is often destroyed in accidents. However, due to the important barrier and physiologic function of the skin, the integrity of the skin is important to the well-being of the individual, and any breach or rupture represents a threat that must be met by the body in order to protect its continued existence.\nApart from injuries on the skin, injuries may also be present in all kinds of tissues (i.e. soft and hard tissues). Injuries on soft tissues including mucosal membranes and/or skin are especially relevant in connection with the present invention.\nHealing of a wound on the skin or on a mucosal membrane undergoes a series of stages that results either in repair or regeneration of the skin or mucosal membrane. In recent years, regeneration and repair have been distinguised as the two types of healing that may occur. Regeneration may be defined as a biological process whereby the architecture and function of lost tissue are completely renewed. Repair, on the other hand, is a biological process whereby continuity of disrupted tissue is restored by new tissues which do not replicate the structure and function of the lost ones.\nThe majority of wounds heal through repair, meaning that the new tissue formed is structurally and chemically unlike the original tissue (scar tissue). In the early stage of the tissue repair, one process which is almost always involved is the formation of a transient connective tissue in the area of tissue injury. This process starts by formation of a new extracellular collagen matrix by fibroblasts. This new extracellular collagen matrix is then the support for a connective tissue during the final healing process. The final healing is, in most tissues, a scar formation containing connective tissue. In tissues which have regenerative properties, such as, e.g., skin and bone, the final healing includes regeneration of the original tissue. This regenerated tissue has frequently also some scar characteristics, e.g. a thickening of a healed bone fracture.\nUnder normally circumstances, the body provides mechanisms for healing injured skin or mucosa in order to restore the integrity of the skin barrier or the mucosa. The repair process for even minor ruptures or wounds may take a period of time extending from hours and days to weeks. However, in ulceration, the healing can be very slow and the wound may persist for an extended period of time, i.e. months or even years.\nThe stages of wound healing normally include inflammation (normally 1-3 days), migration (normally 1-6 days), proliferation (normally 3-24 days ) and maturation (normally 1-12 months). The healing process is a complex and well orchestrated physiological process that involves migration, proliferation and differentiation of a variety of cell types as well as synthesis of matrix components. The healing process may be separated into the following three phases:\ni) Haemostasis and inflammation\nWhen platelets are present outside the circulatory system and exposed to thrombin and collagen, they become activated and they aggregate. Thus, platelets initiate the repair process by aggregating and forming a temporary plug to ensure haemostasis and prevent invasion from bacteria. The activated platelets initiate the coagulation system and release growth factors like platelet-derived growth factor (PDGF) and epidermal growth factors (EGFs) and transforming growth factors (TGFs).\nThe first cells to invade the wound area are neutrophils followed by monocytes which are activated by macrophages.\nThe major role of neutrophils appears to be clearing the wound of or defending the wound against contaminating bacteria and to improve the healing of the wound by removing dead cells and platelets. The infiltration of neutrophils ceases within about the first 48 hours provided that no bacterial contamination is present in the wound. Excess neutrophils are phagocytosed by tissue macrophages recruited from the circulating pool of blood-borne monocytes. Macrophages are believed to be essential for efficient wound healing in that they also are responsible for phagocytosis of pathogenic organisms and a clearing up of tissue debris. Furthermore, they release numerous factors involved in subsequent events of the healing process. The macrophages attract fibroblasts which start the production of collagen.\nii) Granulation Tissue Formation and Re-Epithelization\nWithin 48 hours after wounding, fibroblasts begin to proliferate and migrate into the wounds space from the connective tissue at the wound edge. The fibroblasts produce collagens and glycosaminoglycans and inter alia low oxygen tension at the wound stimulates proliferation of endothelial cells. The endothelial cells give rise to the formation of a new capillary network.\nCollagenases and plasminogen activators are secreted from keratinocytes. If the wound is left undisturbed and well-nourished with oxygen and nutrients, keratinocytes will migrate over the wound. Keratinocytes are believed only to migrate over viable tissue and, accordingly, the keratinocytes migrate into the area below the dead tissue and the crust of the wound.\nThe wound area is further decreased by contraction.\niii) Dermal Remodelling\nAs soon as the re-epithelization is completed the remodelling of the tissue begins. This phase, which lasts for several years, restores the strength to the wounded tissue.\nAll of the above-mentioned healing processes take considerable time. The rate of healing is influenced by the wound\"\"s freedom from infection, the general health of the individual, presence of foreign bodies, etc. Some pathologic conditions like infection, maceration, dehydration, generally poor health and malnutrition can lead to formation of a chronic ulcer such as, e.g., ischemic ulcers.\nUntil at least superficial healing has occurred, the wound remains at risk of continued or new infection. Therefore, the quicker the wound can heal, the sooner the risk is removed.\nThus, any procedure that can influence the rate of wound healing or favourably influence the healing of wounds is of great value.\nFurthermore, as almost all tissue repair processes include the early connective tissue formation, a stimulation of this and the subsequent processes are contemplated to improve tissue healing.\nIn the present context the term xe2x80x9cclinical healingxe2x80x9d is used to denote a situation where not tissue interruption can be visually observed and only discrete signs of inflammation are present such as a light redness or a discretely swollen tissue. In addition, no complaints of pain are present when the organ is relaxed or untouched.\nAs mentioned above, the invention relates to the use of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins as a wound healing agent, i.e. an agent which accelerates, stimulates or promotes healing of dermal or mucosal wounds. Accordingly, an important use is also the use as tissue regeneration and/or repair agents. Furthermore, due to the wound healing effect, enamel matrix, enamel matrix derivatives and/or enamel matrix proteins have pain relief effect.\nTraditionally, dry or wet-to-dry dressings have been most commonly used for wound care. These are gradually being replaced by moist environments using occlusive dressings. To successfully repair or replace a failed body part, the processes of wound healing, fibrosis and microbial invasion must be balanced against each other. Many tools available to ward off infection compromise wound healing. Delayed wound healing or inflammation can exacerbate fibrosis. Moreover, it has previously been suggested that growth factors like epidermal growth factor (EGF), transforming growth factor-xcex1 (TGF-xcex1), platelet derived growth factor (PDGF), fibroblast growth factors (FGFs) including acidic fibroblast growth factor (xcex1-FGF) and basic fibroblast growth factor (xcex2-FGF), transforming growth factor-xcex2 (TGF-xcex2) and insulin like growth factors (IGF-1 and IGF-2) are conductors of the wound healing process and they are frequently cited as promoters of wound healing; however, they can actually promote fibrosis which in turn may impair successful healing. Even though accelerated healing offers the most promise for reducing the risk of infection and the resulting inflammation that can lead to scar formation, therapeutic attempts to accelerate the normal wound healing process have met with relatively little success. This is likely because the repair process involves the concerted involvement of a number of factors, cf. above.\nTo this end, the present inventors have observed that in various cell cultures of fibroblasts (embryonal, dermal, derived from the periodontal ligament, fish or bird), twice as much TGFxcex21 is produced in the cell cultures stimulated with EMDOGAIN(copyright) compared to non-stimulated cultures when assayed by, e.g., ELISA in a sample from the culture medium (vide Example 1 below). The increase is present after 24 hours of culture, but more pronounced on the following days (days 2 and 3). After the second day, also the cell proliferation is increased in cell cultures stimulated with EMDOGAIN(copyright). A similar but less pronounced increase of TGFxcex21 production is observed in human epithelial cells. As TGFxcex21 seems to be of central importance in the epithelisation of surface wounds, these findings support the concept of the present invention.\nIn the oral cavity the use of dressings is common. Such dressings are of the traditional type, e.g. Surgipads to stop bleeding and Coe-Pack periodontal dressing (Coe Laboratories, the GC Group, Chicago, USA) on open wounds, Gaze drenched in antibiotic solution is inserted in tooth extraction alveoli and requires removal after a few days when the healing has started. Rinsing with antiseptics such as chlorhexidine is regularly used after oral surgery. Sometimes general or topical antibiotics are also prescribed.\nIn general specific precautions have to be taken into considerations in connection with treatment of wounds, such as, e.g., sterility considerations, contamination problems, correct application of bandages/dressings etc. which normally require that the treatment/application is performed by well-educated nurses or the like. Thus, wound treatment often becomes a very expensive operation when the wound healing agent is to be applied several times daily. A desired reduction in the costs involved in wound healing treatment is therefore obtainable when the application frequency can be reduced or if the healing processes are improved leading to a reduction in the time period required to heal the wound.\nThe present inventors have now found that enamel matrix, enamel matrix derivatives and/or enamel matrix proteins have wound healing properties. Furthermore, there are indications of that the application of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins lead to improved wound healing. Especially, the inventors have observed that after application of enamel matrix proteins and/or enamel matrix derivatives, the inflammation stage is shortened and the typical signs such as warmth, redness, oedema and pain are less noticeable, and new tissues are formed more rapidly. The observed time for wound healing (e.g. after surgery) is significantly shortened as compared to surgery without the use of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins.\nThe therapeutic and/or prophylactic activity of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins may of course be evidenced by in vivo tests using experimental animals or humans (cf. the experimental section herein). However, an indication of the efficacy and/or activity of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins can be obtained by performing relatively simple in vitro tests such as, e.g., tests involving cell cultures.\nFurthermore, there are several parameters which may be employed in order to evaluate a wound healing effect. These include:\nComputer aided planimetry (evaluation of rate of open wound healing)\nLaser doppler imaging (evaluation of wound perfusion)\nTensiometry (evaluation of wound strength)\nHistopathology/cytology (microscopic evaluation of wound tissues and fluids)\nBiochemistry (HPLC/RIA) (evaluation of various drugs and biochemical components of tissue healing)\nElectrodiagnostics (evaluation of relationship of wound healing and innervation)\nScintigraphy (radionuclide imaging of wound tissue)\nIn connection with treatment of wounds/ulcers, debridement and wound cleansing are of particular importance. It is believed that the cleaning and/or debridement of wounds/ulcers are a prerequisite for the healing process and, furthermore, when wound healing agents are applied such agents have to exert their effect on fresh and vital tissue and not on dead tissue or contaminated tissue. Debridement of necrotic tissue can be performed by at least four different methods: i) sharp debridement, ii) mechanical debridement, iii) enzymatic debridement, and iv) autolytic debridement.\nTherefore, the present invention relates also to the use of a debridement method in combination with the use of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins for the healing or prevention of wounds. Such combination therapy involves the following two steps, namely i) a debridement method and ii) application of an enamel matrix, enamel matrix derivatives and/or enamel matrix proteins and the two steps may be carried out as many times as desired and in any suitable order.\nWhen the wound has been subjected to debridement, the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins may be applied either directly on or into the wound or it can be applied in the form of any suitable pharmaceutical composition such as, e.g., a dry or moist, clean dressing into which the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins has been incorporated. The enamel matrix, enamel matrix derivatives and/or enamel matrix proteins may of course also be applied in connection with cleansing of the wound.\nAs will be discussed later, the enamel matrix, enamel matrix derivatives and/or the enamel matrix proteins may be used as such or they may be used in a suitable preparation or pharmaceutical composition.\nInfection-Decreasing Effect\nIn a further aspect of the present invention, the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins are used as therapeutic or prophylactic agents having an anti-microbial effect. The enamel matrix, enamel matrix derivatives and/or enamel matrix proteins exhibit infection-decreasing properties.\nIn the present context the term infection-decreasing effect relates to a treating or preventive effect by the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins on an infection in a tissue of an individual when the tissue or the individual is treated with the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins.\nThe term infection relates to the invasion and multiplication of microorganisms in body tissues or accumulation on the tissues, which may be clinically inapparent or result in local cellular injury due to competitive metabolism, enzymes, toxins, intracellular replication or antigen-antibody response.\nIn accordance with the present invention, the infection to be prevented and/or treated may be caused by a microorganism. The microorganisms of interest according to the present invention include bacteria, viruses, yeast, molds, protozoa and rickettsiae.\nIn the present context the term xe2x80x9canti-bacterial effectxe2x80x9d means that the growth of bacteria is suppressed or the bacteria are destroyed. The term is not limited to certain bacteria but encompasses in general any bacteria. However, the invention is focused on i) pathogenic bacteria which cause diseases in mammals including humans and/or ii) bacteria which normally are present in a mammal body and which under certain conditions may cause unwanted conditions in the body.\nAccordingly, the invention relates to the use of an active enamel substance for the prevention of or treatment of bacterial growth on a body surface such as the skin, a mucosal surface or a nail or a tooth surface.\nGeneral and specific description of the bacterial conditions to be counteracted\nThe enamel matrix, enamel matrix derivatives and/or enamel matrix proteins may be used for the treatment of an infection caused by bacteria together with or without the presence of an antimicrobial. Gram negative bacteria to be treated with the active enamel substance could be cocci, such as Neisseria (e.g. N. meningitis, N. gonorrhoeae), and Acinetobacter or rods, such as Bacteriodes (e.g. B. fragilis), Bordetella (e.g. B. pertussis, B. parapertussis), Brucella (e.g. B. melitentis, B. abortus Bang, B suis), Campylobacter (e.g. C jejuni, C. coli, C. fetus), Citrobacter, Enterobacter, Escherichia (e.g. E. coli), Haemophilus (e.g. H. influenzae, H. parainfluenzae), Klebsiella (e.g. K. pneumoniae), Legionella (e.g. L. pneumophila), Pasteurella (e.g. P. yersinia, P. multocida), Proteus (e.g. P. mirabilis, P. vulgaris), Pseudomonas (e.g. P. aeruginosa, P. pseudomallei, P. mallei), Salmonella (e.g. S. enteritidis, S. infantitis S. Dublin S. typhi, S. paratyphi, S. schottmxc3xclleri, S. choleraesuis, S. typhimurium, or any of the 2,500 other serotypes), Serratia (e.g. S. marscences, S. liquifaciens), Shigella (e.g. S. sonnei, S. flexneri, S. dysenteriae, S. boydii), Vibrio (e.g. V. cholerae, V. el tor), and Yersinia (e.g. Y. enterocolitica, Y. pseudotuberculosis, Y. pestis). Gram positive bacteria to be treated with the active enamel substance could be cocci, such as Streptococcus (e.g. S. pneumoniae, S. viridans, S. faecalis, S. pyogenes), Staphylococcus (e.g. S. aureus, S. epidermidis, S. saprophyticus, S. albus), and rods, such as Actinomyces (e.g. A. israelli), Bacillus (e.g. B. cereus, B. subtilis, B. anthracis), Clostridium (e.g. C. botulinum, C. tetani, C. perfringens, C. difficile), Corynebacterium (e.g. C. diphtheriae), Listeria, and Providencia. Other bacteria causing infection include Propinobacterium acne and Pityosporon ovale. \nThe enamel matrix, enamel matrix derivatives and/or enamel matrix proteins may also be used for the treatment of an infection caused by a spirochete such as, e.g., Borrelia, Leptospira, Treponema or Pseudomonas.\nAn antimicrobial to be used in combination with the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins could be an antimicrobial that has an anti-microbial action through inhibition of cell wall synthesis, such as xcex2-lactams and vancomycin, preferably penicillins, such as amdinocillin, ampicillin, amoxicillin, azlocillin, bacampicillin, benzathine pinicillin G, carbenicillin, cloxacillin, cyclacillin, dicloxacillin, methicillin, mezlocillin, nafcillin, oxacillin, penicillin G, penicillin V, piperacillin, and ticarcillin;\ncephalosporins, such as the first generation drugs cefadroxil, cefazolin, cephalexin, cephalothin, cephapirin, and cephradine, the second generation drugs cefaclor, cefamandole, cefonicid, ceforanide, cefoxitin, and cefuroxime, or the third generation cephalosporins cefoperazone, cefotaxime, cefotetan, ceftazidime, ceftizoxime, ceftriaxone, and moxalactam; carbapenems such as imipenem; or monobactams such as aztreonam.\nOther antimicrobial drugs with action through inhibition of protein synthesis, such as chloramphenicol; other tetracyclines preferably demeclocycline, doxycycline, methacycline, minocycline, and oxytetracycline; aminoglycosides such as amikacin, gentamicin, kanamycin, neomycin, netilmicin, paromomycin, spectinomycin, streptomycin, and tobramycin; polymyxins such as colistin, colistimathate, and polymyxin B, and erythromycins and lincomycins;\nantimicrobials with action through inhibition of nucleic acid synthesis in particular sulfonamides such as sulfacytine, sulfadiazine, sulfisoxazole, sulfamethoxazole, sulfamethizole, and sulfapyridine; trimethoprim, quinolones, novobiocin, pyrimethanmine, and rifampin.\nIn a specific embodiment of the invention, the infection is present in the oral cavity and the infection may be a bacterial condition.\nOral bacteria to be contact inhibited or otherwise combated. Examples (not conditions) include\nbacteria causing caries, e.g. Streptococcus mutans, Lactobacillus spp.\nbacteria causing periodontal disease e.g. Actinobacillus actinomycetemcomitans, Porphyromonas gingivalis, Prevotella intermedia, Peptostreptococcus micros, Campylobacter (Fusobacteria, Staphylococci), B. forsythus \nbacteria causing alveolitis etc., e.g. Staphylococcus, Actinomyces and Bacillus\nbacteria causing periapical lesions, e.g. Spirochetes and all above\nAnti-Inflammatory Effect\nThe present invention also relates to the uses of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins as therapeutic or porphylactive agents having an anti-inflammatory effect.\nSeveral drugs are employed to suppress the manifestations of inflammation, including the ardrenocorticosteroids, the large group comprising the so called non-steroid anti-inflammatory drugs or NSAIDs, and drugs such as immunosuppressive agents. Adrenocorticosteroids, and especially glucocorticoids, have potent anti-inflammatory effects when used in pharmacological doses. They specifically inhibit the early vascular phase of the inflammatory process by decreasing the vascular permeability and thereby granulocyte migration. Glucocorticoids also interfere with late inflammatory and reparative processes, in that they inhibit the proliferation of mesenchymal cells and the production of extracellular macromolecules, including proteoglycanes and collagen. It has been shown experimentally that glucocorticoids inhibit, for example, macrophage function, production of humoral antibodies, cellular immunity, and possibly the release of lysosomal enzymes.\nThe severity of tissue damage may depend on the antigen/antibody reaction of the organism as well as the degree of retention of inflammatory products in the affected area. Accumulation of mediators of local inflammation accelerates the process. In most cases the process is slow, with immunoinfiltration of the tissue and formation of granulation tissue which contains inflammatory cells.\nIn the present context the term xe2x80x9canti-inflammatory effectxe2x80x9d denotes a counteracting or suppression of inflammation.\nGeneral and specific description of the kind of inflammatory conditions to be treated\nThe inflammatory condition to be treated in accordance with the present invention may of course be any inflammatory condition in/on any part of the body or any inflammatory condition present in soft or hard tissue. In one embodiment of the invention the inflammatory condition is present in the oral cavity. Examples of conditions in the oral cavity are alveolitis, cheilitis, bone necrosis (after trauma), fractures.\nIn another embodiment of the invention, the inflammatory condition is present in a bone donor site. In a third embodiment of the invention, the inflammatory condition is present in a joint cavity. Examples of such inflammatory conditions are rheumatoid arthritis and related conditions.\nAnti-Bacterial Versus Anti-Inflammatory\nIn contrast to many currently used antibiotic agents, enamel matrix proteins will not compromise wound healing and the rapid wound healing in turn does not leave room for chronic or long lasting inflammation processes to develop. Also, the reorganisation of proper tissues, such as described after application of enamel matrix derivatives onto periodontal defects, is clearly favoured by a rapid wound healing without bacteria or inflammatory reactions.\nThe application of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins leads to rapid wound healing of surgical incisions, possibly by creating a surface which in contact with bacteria inhibit their growth but a the same time enhances fibroblast migration and collagen synthesis. If the inflammatory stage is shortened, the typical signs such as warmth, redness, oedema and pain are less noticeable.\nEnamel Matrix, Enamel Matrix Derivatives and Enamel Matrix Proteins\nEnamel matrix is a precursor to enamel and may be obtained from any relevant natural source, i.e. a mammal in which teeth are under development. A suitable source is developing teeth from slaughtered animals such as, e.g., calves, pigs or lambs. Another source is for example fish skin.\nEnamel matrix can be prepared from developing teeth as described previously (EP-B-0 337 967 and EP-B-0 263 086). The enamel matrix is scraped off and enamel matrix derivatives are prepared, e.g. by extraction with aqueous solution such as a buffer, a dilute acid or base or a water/solvent mixture, followed by size exclusion, desalting or other purification steps, optionally followed by freeze-drying. Enzymes may be deactivated by treatment with heat or solvents, in which case the derivatives may be stored in liquid form without freeze-drying.\nIn the present context, enamel matrix derivatives are derivatives of enamel matrix which include one or several of enamel matrix proteins or parts of such proteins, produced naturally by alternate splicing or processing, or by either enzymatic or chemical cleavage of a natural length protein, or by synthesis of polypeptides in vitro or in vivo (recombinant DNA methods or cultivation of diploid cells). Enamel matrix protein derivatives also include enamel matrix related polypeptides or proteins. The polypeptides or proteins may be bound to a suitable biodegradable carrier molecule, such as polyamino acids or polysaccharides, or combinations thereof. Furthermore, the term enamel matrix derivatives also encompasses synthetic analogous substances.\nProteins are biological macromolecules constituted by amino acid residues linked together by peptide bonds. Proteins, as linear polymers of amino acids, are also called polypeptides. Typically, proteins have 50-800 amino acid residues and hence have molecular weights in the range of from about 6,000 to about several hundred thousand Daltons or more. Small proteins are called peptides or oligopeptides.\nEnamel matrix proteins are proteins which normally are present in enamel matrix, i.e. the precursor for enamel (Ten Cate: Oral Histology, 1994; Robinson: Eur. J. Oral Science, January 1998, 106 Suppl. 1:282-91), or proteins which can be obtained by cleavage of such proteins. In general such proteins have a molecular weight below 120,000 daltons and include amelogenins, non-amelogenins, proline-rich non-amelogenins, amelins (ameloblastin, sheathlin) and tuftelins.\nExamples of proteins for use according to the invention are amelogenins, proline-rich non-amelogenins, tuftelin, tuft proteins, serum proteins, salivary proteins, amelin, ameloblastin, sheathlin, and derivatives thereof, and mixtures thereof. A preparation containing an active enamel substance for use according to the invention may also contain at least two of the aforementioned proteinaceous substances. A commercial product comprising amelogenins and possibly other enamel matrix proteins is marketed as EMDOGAIN(copyright) (Biora AB).\nIn general, the major proteins of an enamel matrix are known as amelogenins. They constitute about 90% w/w of the matrix proteins. The remaining 10% w/w includes proline-rich non-amelogenins, tuftelin, tuft proteins, serum proteins and at least one salivary protein; however, other proteins may also be present such as, e.g., amelin (ameloblastin, sheathlin) which have been identified in association with enamel matrix. Furthermore, the various proteins may be synthesized and/or processed in several different sizes (i.e. different molecular weights). Thus, the dominating proteins in enamel matrix, amelogenins, have been found to exist in several different sizes which together form supramolecular aggregates. They are markedly hydrophobic substances which under physiologically conditions form aggregates. They may carry or be carriers for other proteins or peptides.\nOther protein substances are also contemplated to be suitable for use according to the present invention. Examples include proteins such as proline-rich proteins and polyproline. Other examples of substances which are contemplated to be suitable for use according to the present invention are aggregates of such proteins, of enamel matrix derivatives and/or of enamel matrix proteins as well as metabolites of enamel matrix, enamel matrix derivatives and enamel matrix proteins. The metabolites may be of any size ranging from the size of proteins to that of short peptides.\nAs mentioned above, the protein, polypeptides or peptides for use according to the invention typically have a molecular weight of at the most about 120 kDa such as, e.g., at the most 100 kDa, 90 kDa, 80 kDa, 70 kDa or 60 kDa as determined by SDS Page electrophoresis.\nThe proteins for use according to the invention are normally presented in the form of a preparation, wherein the protein content of the active enamel substance in the preparation is in a range of from about 0.05% w/w to 100% w/w such as, e.g., about 5-99% w/w, about 10-95% w/w, about 15-90% w/w, about 20-90% w/w, about 30-90% w/w, about 40-85% w/w, about 50-80% w/w, about 60-70% w/w, about 70-90% w/w, or about 80-90% w/w.\nA preparation of an active enamel substance for use according to the invention may also contain a mixture of active enamel substances with different molecular weights.\nThe proteins of an enamel matrix can be divided into a high molecular weight part and a low molecular weight part, and it has been found that a well-defined fraction of enamel matrix proteins possesses valuable properties with respect to treatment of periodontal defects (i.e. periodontal wounds). This fraction contains acetic acid extractable proteins generally referred to as amelogenins and constitutes the low molecular weight part of an enamel matrix (cf. EP-B-0 337 967 and EP-B-0 263 086).\nAs discussed above the low molecular weight part of an enamel matrix has a suitable activity of inducing binding between hard tissues in periodontal defects. In the present context, however, the active proteins are not restricted to the low molecular weight part of an enamel matrix. At present, preferred proteins include enamel matrix proteins such as amelogenin, amelin, tuftelin, etc. with molecular weights (as measured in vitro with SDS-PAGE) below about 60,000 daltons but proteins having a molecular weight above 60,000 daltons have also promising properties as candidates for wound healing, anti-bacterial and/or anti-inflammatory agents.\nAccordingly, it is contemplated that the active enamel substance for use according to the invention has a molecular weight of up to about 40,000 such as, e.g. a molecular weight of between about 5,000 and about 25,000.\nWithin the scope of the present invention are also peptides as described in WO 97/02730, i.e. peptides which comprise at least one sequence element selected from the group consisting of the tetrapeptides DGEA (Asp-Gly-Glu-Ala), VTKG (Val-Thr-Lys-Gly), EKGE (Glu-Lys-Gly-Glu) and DKGE (Asp-Lys-Gly-Glu) and which further comprise an amino acid sequence from which a consecutive string of 20 amino acids is identical to a degree of at least 80% with a string of amino acids having the same length selected from the group consisting of the amino acid sequence shown in SEQ ID NO: 1 and a sequence consisting of amino acids 1 to 103 of SEQ ID NO: 1 and amino acids 6 to 324 of SEQ ID NO:2.\nBy the term xe2x80x9csequence identityxe2x80x9d is meant the identity in sequence of amino acids in the match with respect to identity and position of the amino acids of the peptides. A gap is counted as non-identity for one or more amino acids as appropriate.\nSuch peptides may comprise from 6 to 300 amino acids, e.g. at least 20 amino acids, at least 30 amino acids, such as at least 60 amino acids, at least 90 amino acids, at least 120 amino acids, at least 150 amino acids or at least 200 amino acids.\nA method for the isolation of enamel matrix proteins involves extraction of the proteins and removal of calcium and phosphate ions from solubilized hydroxyapatite by a suitable method, e.g. gel filtration, dialysis of ultrafiltration (see e.g. Janson, J-C and Rydxc3xa9n, L. (Eds.), Protein purification, VCH Publishers 1989 and Harris, ELV and Angal, S., Protein purification methodsxe2x80x94A practical approach, IRL Press, Oxford 1990).\nA typical lyophilized protein preparation may mainly or exclusively up to 70-90% contain amelogenins with a molecular weight (MW) between 40,000 and 5,000 daltons, the 10-30% being made up of smaller peptides, salts and residual water. The main protein bands are at 20 kDa and around 5 kDa.\nBy separating the proteins, e.g. by precipitation, ion-exchange chromatography, preparative electrophoresis, gel permeation chromatography, reversed phase chromatography of affinity chromatography, the different molecular weight amelogenins can be purified.\nThe combination of molecular weight amelogenins may be varied, from a dominating 20 kDa compound to an aggregate of amelogenins with many different molecular weights between 40 and 5 kDa, and to a dominating 5 kDa compound. Other enamel matrix proteins such as amelin, tuftelin or proteolytic enzymes normally found in enamel matrix, can be added and carried by the amelogenin aggregate.\nAs an alternative source of the enamel matrix derivatives or protein one may also use generally applicable synthetic routes well-known for a person skilled in the art or use cultivated cells or bacteria modified by recombinant DNA-techniques (see, e.g., Sambrook, J. et al.: Molecular Cloning, Cold Spring Harbor Laboratory Press, 1989).\nPhysico-chemical properties of enamel matrix, enamel matrix derivatives and enamel matrix proteins\nIn general the enamel matrix, enamel matrix derivatives and enamel matrix proteins are hydrophobic substances. i.e. less soluble in water especially at increased temperatures. In general, these proteins are soluble at non-physiological pH values and at a low temperature such as about 4-20xc2x0 C. while they will aggregate and precipitate at body temperature (35-37xc2x0 C.) and neutral pH.\nThe enamel matrix, enamel matrix derivatives and/or enamel matrix proteins for use according to the invention also include an active enamel substance, wherein at least a part of the active enamel substance is in the form of aggregates or after application in vivo is capable of forming aggregates. The particle size of the aggregates is in a range of from about 20 nm to about 1 xcexcm.\nIt is contemplated that the solubility properties of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins are of importance in connection with the prophylactive and therapeutic activity of the substances. When a composition containing the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins (in the following also denoted xe2x80x9cactive enamel substancexe2x80x9d as a common term) is administered to e.g. a human, the proteinaceous substances will precipitate due to the pH normally prevailing under physiological conditions. Thus, a layer of enamel matrix, enamel matrix derivatives and/or enamel proteins is formed at the application site and this layer (which also may be a molecular layer in those cases where aggregates have been formed) is difficult to rinse off under physiological conditions. Furthermore, due to the substances bioadhesive properties (see below) the precipitated layer is firmly bound to the tissue also at the margin between the precipitated layer and the tissue. The proteinaceous laeyr thus covers the tissue onto which the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins or compositions thereof have been applied and the active enamel substances are maintained in situ for a prolonged period of time, i.e. it is not necessary to administer the active enamel substances(s) with short intervals. Furthermore, the layer formed in situ can almost be compared to an occlusive dressing, i.e., the layer formed protects the tissue onto which the layer is formed from the surroundings. In the case of a wound tissue, an infected tissue of an inflamed tissue such a layer protects the tissue from further contamination from microorganisms present in the surroundings. Furthermore, the proteinaceous layer may exert its effect by direct contact with the tissue or with microorganisms present in/on/at the tissue.\nIn order to enable a proteinaceous layer to be formed in situ after application it may be advantageous to incorporate a suitable buffer substance in a pharmaceutical or cosmetic composition of the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins; the purpose of such a buffer substance could be to avoid the dissolution of the active enamel substance at the application site.\nThe enamel matrix, enamel matrix derivatives and enamel matrix proteins have also been observed (by the present inventors) to posses bioadhesives properties, i.e. they have an ability to adhere to skin or mucosal surfaces. These properties are most valuable in connection with a therapeutic and/or prophylactic treatment at least for the following reasons:\nthe prophylactically and/or therapeutically active substance(s) can be maintained at the application site for a prolonged period of time (i.e. i) the administration frequency can be reduced, ii) a controlled release effect of the active substance is obtainable and/or iii) a local treatment at the application site is improved)\nthe substances may in themselves be suitable as vehicles for other prophylactically or therapeutically active substances because a vehicle containing enamel matrix, enamel matrix derivatives and/or enamel matrix proteins can be formulated as a bioadhesive vehicle (i.e. a novel bioadhesive drug delivery system based on the bioadhesive properties of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins).\nTheories with Respect to Mechanism of Action\nEnamel matrix is an example of extracellular protein matrix which adheres to mineral surfaces as well as to proteinaceous surfaces. At physiological pH and temperature the proteins form an insoluble supra-molecular aggregate (Fincham et al. in J. Struct. Biol. 1994 March-April; 112(2):103-9 and in J. Struct. Biol. 1995 July-August; 115(1):50-9), which is gradually degraded by proteolytic enzymes (occurs both in vivo and in vitro provided that the proteases have not been subjected to inactivation).\nThe recent observation that enamel matrix is formed and temporarily present during root and root cementum formation can explain how application of enamel matrix, enamel matrix derivatives and/or enamel matrix proteins promotes the regeneration of periodonetal tissue. However, the observation underlying the present invention that enamel matrix, enamel matrix derivatives and/or enamel matrix proteins also have a positive effect on healing of soft tissue defects like wound healing is very surprising. The same applies to the observations with respect to anti-infectious and anti-inflammatory effect.\nIn many species, remnants of enamel matrix are found in the newly mineralized crown when a tooth is erupting into the oral cavity. It might be argued that a new tooth would be very vulnerable to bacterial attack from common oral bacteria unless it had a natural protection during this initial phase.\nApplication of insolubilising enamel matrix, enamel matrix derivatives and/or enamel matrix proteins with suitable anti-bacterial and/or anti-inflammatory properties onto a wounded surface will enhance and improve healing.\nAs demonstrated in the experimental section herein the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins or protein aggregates hinder bacterial growth by contact inhibition, while exposed cells apparently react upon the enamel matrix as a normal environment which suppresses inflammatory responses.\nIn accordance with the present invention, an enamel matrix, enamel matrix derivative and/or enamel matrix protein may be used for curative purposes as well as for preventive purposes. Furthermore, an enamel matrix, enamel matrix derivative and/or enamel matrix protein may be used together with other active drug substances such as, e.g. anti-bacterial, anti-inflammatory, antiviral, antifungal substances or in combination with growth factors such as, e.g., TGFxcex2, PDGF, IGF, FGF, keratinocyte growth factor or peptide analogues thereof (it is believed that EGF promotes healing by enhancing migration and cell division of epithelial cells; furthermore, EGF increases fibroblast numbers in wound resulting in a greater collagen production). Enzymesxe2x80x94either inherently present in the enamel matrix or preparation thereof or addedxe2x80x94may also be used in combination with an enamel matrix, enamel matrix derivative and/or enamel matrix protein, especially proteases.\nA preparation of the active enamel substances is normally formulated as a pharmaceutical or cosmetic composition. Such as composition may of course consist of the proteinaceous preparation or it may further comprise a pharmaceutically or cosmetically acceptable excipient. Especially suitable excipients for use in pharmaceutic or cosmetic compositions are propylene glycol alginate, or hyaluronic acid or salts or derivatives thereof.\nPharmaceutical and/or Cosmetic Compositions\nIn the following examples of suitable compositions containing the active enamel substance(s) are given. Depending on the use of the active enamel substance(s), a composition may be a pharmaceutical or a cosmetic composition. In the following the term xe2x80x9cpharmaceutical compositionxe2x80x9d is also intended to embrace cosmetic compositions as well as compositions belonging to the so-called grey area between pharmaceuticals and cosmetics, namely cosmeceuticals.\nFor the administration to an individual (an animal or a human) the enamel matrix, enamel matrix derivatives and/or enamel matrix proteins (in the following also denoted xe2x80x9cactive enamel substancexe2x80x9d) and/or a preparation thereof are preferably formulated into a pharmaceutical composition containing the active enamel substance and, optionally, one or more pharmaceutically acceptable excipients.\nThe compositions may be in form of, e.g., solid, semi-solid or fluid compositions such as, e.g.,\nbioabsorbable patches, drenches, dressings, hydrogel dressings, hydrocolloid dressings, films, foams, sheets, bandages, plasters, delivery devices, implants,\npowders, granules, granulates, capsules, agarose or chitosan beads, tablets, pills, pellets, microcapsules, microspheres, nanoparticles,\nsprays, aerosols, inhalation devices,\ngels, hydrogels, pastes, ointments, creams, soaps, suppositories, vagitories, tooth paste,\nsolutions, dispersions, suspensions, emulsions, mixtures, lotions, mouthwash, shampoos, enemas,\nkits containing e.g. two separate containers, wherein the first one of the containers contains the active enamel substance optionally admixed with other active drug substance(s) and/or pharmaceutically acceptable excipients and the second container containing a suitable medium intended to be added to the first container before use in order to obtain a ready-to-use composition;\nand in other suitable forms such as, e.g., implants or coating of implants or in a form suitable for use in connection with implantation or transplantation.\nCompositions for application to the skin or to the mucosa are considered most important in connection with the present invention. Thus, a composition comprising the active enamel substance to be administered may be adapted for administration by any suitable route, for example by topical (dermal), oral, buccal, nasal, aural, rectal or vaginal administration, or by administration to a body cavity such as, e.g., a tooth root or a tooth root canal. Furthermore, a composition may be adapted to administration in connection with surgery, e.g. in connection with incision within the body in order to promote healing of internal wounds and soft tissue damages.\nAs mentioned above, a composition of the active enamel substance(s) may be suitable for use during surgery, e.g. for local application (e.g. in the oral cavity) in the form of a gel, film or dry pellet, or as a rinsing solution or treatment with a paste or cream on tissue or surfaces to prevent bacterial attack. In connection with surgery or implantation in the area of the tooth root canal, a paste for cavity sealing can be employed.\nThe compositions may be formulated according to conventional pharmaceutical practice, see, e.g., xe2x80x9cRemington\"\"s Pharmaceutical Sciencesxe2x80x9d and xe2x80x9cEncyclopedia of Pharmaceutical Technologyxe2x80x9d, edited by Swarbrick, J. and J. C. Boylan, Marcel Dekker, Inc., New York, 1988.\nAs mentioned above, the application of a composition comprising an active enamel substance is intended for skin or mucosa. Other applications may of course also be relevant such as, e.g., application on dentures, protheses, implants, and application to body cavities such as the oral, nasal and vaginal cavity. The mucosa is preferably selected from oral, buccal, nasal, aural, rectal and vaginal mucosa. Furthermore, the application may be directly on or onto a wound or other soft tissue injuries.\nFurthermore, application within the dental/odontologic area is also of great importance. Relevant examples are application to periodontal (dental) pockets, to gingiva or to gingival wounds or other wounds located in the oral cavity, or in connection with oral surgery.\nIt is further anticipated that, due to the antibacterial properties of the active enamel substance described herein, it may advantageously be applied to teeth or tooth roots for the prevention of caries and/or plaque. To support this use, it has been shown that teeth which are imperfectly developed (amelogenesis imperfecta) and consequently contain large amounts of amelogenins are remarkably caries resistant.\nA pharmaceutical composition comprising an active enamel substance serves as a drug delivery system. In the present context the term xe2x80x9cdrug delivery systemxe2x80x9d denotes a pharmaceutical composition (a pharmaceutical formulation or a dosage form) which upon administration presents the active substance to the body of a human or an animal. Thus, the term xe2x80x9cdrug delivery systemxe2x80x9d embraces plain pharmaceutical compositions such as, e.g., creams, ointments, liquids, powders, tablets, etc. as well as more sophisticated formulations such as sprays, plasters, bandages, dressings, devices, etc.\nApart from the active enamel substance, a pharmaceutical composition for use according to the invention may comprise pharmaceutically or cosmetically acceptable excipients.\nA pharmaceutically or cosmetically acceptable excipient is a substance which is substantially harmless to the individual to which the composition is to be administered. Such an excipient normally fulfills the requirements given by the national health authorities. Official pharmacopoeias such as e.g. the British Pharmacopoeia, the United States of America Pharmacopoeia and The European Pharmacopoeia set standards for pharmaceutically acceptable excipients.\nWhether a pharmaceutically acceptable excipient is suitable for use in a pharmaceutical composition is generally dependent on which kind of dosage form is chosen for use for a particular kind of wound. In the following are given examples of suitable pharmaceutically acceptable excipients for use in different kinds of compositions for use according to the invention.\nIn the following is given a review on relevant pharmaceutical compositions for use according to the invention. The review is based on the particular route of administration. However, it is appreciated that in those cases where a pharmaceutically acceptable excipient may be employed in different dosage forms or compositions, the application of a particular pharmaceutically acceptable excipient is not limited to a particular dosage form or of a particular function of the excipient.\nThe choice of pharmaceutically acceptable excipient(s) in a composition for use according to the invention and the optimum concentration thereof cannot generally be predicted and must be determined on the basis of an experimental evaluation of the final composition. However, a person skilled in the art of pharmaceutical formulation can find guidance in e.g., xe2x80x9cRemington\"\"s Pharmaceutical Sciencesxe2x80x9d, 18th Edition, Mack Publishing Company, Easton, 1990.\nTopical compositions\nFor application to the mucosa or the skin, the compositions for use according to the invention may contain conventionally non-toxic pharmaceutically acceptable carriers and excipients including microspheres and liposomes.\nThe compositions for use according to the invention include all kinds of solid, semi-solid and fluid compositions. Compositions of particular relevance are e.g. pastes, ointments, hydrophilic ointments, creams, gels, hydrogels, solutions, emulsions, suspensions, lotions, liniments, shampoos, jellies, soaps, sticks, sprays, powders, films, foams, pads, sponges (e.g. collagen sponges), pads, dressings (such as, e.g., absorbent wound dressings), drenches, bandages, plasters and transdermal delivery systems.\nThe pharmaceutically acceptable excipients may include solvents, buffering agents, preservatives, humectants, chelating agents, antioxidants, stabilizers, emulsifying agents, suspending agents, gel-forming agents, ointment bases, penetration enhancers, perfumes, and skin protective agents.\nExamples of solvents are e.g. water, alcohols, vegetable or marine oils (e.g. edible oils like almond oil, castor oil, cacao butter, coconut oil, corn oil, cottonseed oil, linseed oil, olive oil, palm oil, peanut oil, poppyseed oil, rapeseed oil, sesame oil, soybean oil, sunflower oil, and teaseed oil), mineral oils, fatty oils, liquid paraffin, polyethylene glycols, propylene glycols, glycerol, liquid polyalkylsiloxanes, and mixtures thereof.\nExamples of buffering agents are e.g. citric acid, acetic acid, tartaric acid, lactic acid, hydrogenphosphoric acid, diethylamine etc.\nSuitable examples of preservatives for use in compositions are parabens, such as methyl, ethyl, propyl p-hydroxybenzoate, butylparaben, isobutylparaben, isopropylparaben, potassium sorbate, sorbic acid, benzoic acid, methyl benzoate, phenoxyethanol, bronopol, bronidox, MDM hydantoin, iodopropynyl butylcarbamate, EDTA, benzalconium chloride, and benzylalcohol, or mixtures of preservatives.\nExamples of humectants are glycerin, propylene glycol, sorbitol, lactic acid, urea, and mixtures thereof.\nExamples of chelating agents are sodium EDTA and citric acid.\nExamples of antioxidants are butylated hydroxy anisole (BHA), ascorbic acid and derivatives thereof, tocopherol and derivatives thereof, cysteine, and mixtures thereof.\nExamples of emulsifying agents are naturally occurring gums, e.g. gum acacia or gum tragacanth; naturally occurring phosphatides, e.g. soybean lecithin; sorbitan mono-oleate derivatives; wool fats; wool alcohols; sorbitan esters; monoglycerides; fatty alcohols;, fatty acid esters (e.g. triglycerides of fatty acids); and mixtures thereof.\nExamples of suspending agents are e.g. celluloses and cellulose derivatives such as, e.g., carboxymethyl cellulose, hydroxyethylcellulose, hydroxypropylcellulose, hydroxypropylmethylcellulose, carraghenan, acacia gum, arabic gum, tragacenth, and mixtures thereof.\nExamples of gel bases, viscosity-increasing agents or components which are able to take up exudate from a wound are: liquid paraffin, polyethylene, fatty oils, colloidal silica or aluminium, zinc soaps, glycerol, propylene glycol, tragacanth, carboxyvinyl polymers, magnesium-aluminium silicates, Carbopol(copyright), hydrophilic polymers such as, e.g. starch or cellulose derivatives such as, e.g., carboxymethylcellulose, hydroxyethylcellulose and other cellulose derivatives, water-swellable hydrocolloids, carragenans, hyaluronates (e.g. hyaluronate gel optionally containing sodium chloride), and alginates including propylene glycol aginate.\nExamples of ointment bases are e.g. beeswax, paraffin, cetanol, cetyl palmitate, vegetable oils, sorbitan esters of fatty acids (Span), polyethylene glycols, and condensation products between sorbitan esters of fatty acids and ethylene oxide, e.g. polyoxyethylene sorbitan monooleate (Tween).\nExamples of hydrophobic or water-emulsifying ointment bases are paraffins, vegetable oils, animal fats, synthetic glycerides, waxes, lanolin, and liquid polyalkylsiloxanes.\nExamples of hydrophilic ointment bases are solid macrogols (polyethylene glycols).\nOther examples of ointment bases are triethanolamine soaps, sulphated fatty alcohol and polysorbates.\nExamples of powder components are: alginate, collagen, lactose, powder which is able to form a gel when applied to a wound (absorbs liquid/wound exudate). Normally, powders intended for application on large open wounds must be sterile and the particles present must be micronized.\nExamples of other excipients are polymers such as carmelose, sodium carmelose, hydroxypropylmethylcellulose, hydroxyethylcellulose, hydroxypropylcellulose, pectin, xanthan gum, locust bean gum, acacia gum, gelatin, carbomer, emulsifiers like vitamin E, glyceryl stearates, cetanyl glucoside, collagen, carrageenan, hyaluronates and alginates and kitosans.\nDressings and/or bandages are also important delivery systems for an active enamel substance. When dressings are used as dosage form, the active enamel substance may be admixed with the other material/ingredients before or during the manufacture of the dressing or, the active enamel substance may in some way be coated onto the dressing e.g. by dipping the dressing in a solution or dispersion of the active enamel substance or by spraying a solution or dispersion of the active enamel substance onto the dressing. Alternatively, the active enamel substance may be applied in the form of a powder to the dressing. Dressings may be in the form of absorbent wound dressings for application to exuding wounds. Dressings may also be in the form of hydrogel dressings (e.g. cross-linked polymers such as, e.g. Intrasite(copyright) which contains carboxymethylcellulose, propylene glycol or polysaccharide, disaccharide and proteins) or in the form of occlusive dressings such as, e.g., alginates, chitosan, hydrophilic polyurethane film, collagen sheets, plates, powders, foams, or sponges, foams (e.g. polyurethane or silicone), hydrocolloids (e.g. carboxymethylcellulose, CMC), collagen and hyaluronic acid-based dressings including combinations thereof.\nAlginate, chitosan and hydrocolloid dressings take up wound exudate when placed on a wound. When doing so they produce an aqueous gel on the surface of the wound and this gel is believed to be beneficial for the healing of the wound due to the retaining of moisture in the wound.\nIt is also envisaged that the active enamel substance may be incorporated in a tissue adhesive also comprising, e.g. fibrinogen and thrombin and optionally Factor XIII or another plasma coagulation factor to provide hemostasis. The tissue adhesive may either be prepared as a premix of the active enamel substance, fibrinogen and optionally Factor XIII, thrombin being added to the premix immediately before the tissue adhesive is applied on the wound. Alternatively, the premix of fibrinogen and active enamel substance and optionally Factor XIII may be applied on the wound before application of thrombin. In situ, the thrombin converts fibrinogen to fibrin thereby reproducing the coagulation process occurring naturally in wound healing. The presence of the active enamel substance in the tissue adhesive may serve to accelerate the wound healing process as discussed above. A commercial product suitable for inclusion of the active enamel substance is Tisseel(copyright), a two-component fibrin sealant produced by immuno, AG, Vienna, Austria.\nIn a toothpaste or mouthwash formulation or other formulation for application to teeth or tooth roots, the active enamel substance may either be present in a dissolved state in a vehicle of slightly acid pH or as a dispersion in a vehicle of neutral pH. It is anticipated that in use the active enamel substance may form a protective layer on the surface of the teeth, thereby preventing the attachment of caries producing bacteria (cf. Example 4 below). In such dental care preparations, the active enamel substance may be formulated together with one or more other compounds which have a caries preventive effect, notably fluorine or another trace element such as vanadium or molybdenum. At neutral pH, the trace element is believed to be bound to (e.g. by ion bonds) or embedded in the active enamel substance from which it is released to exert its caries preventive effect when the active enamel substance is dissolved at a pH of about 5.5 or less, e.g. due to acid production by caries producing bacteria.\nThe compositions mentioned above for topical administration are most suitably for application directly to wounds or they may be suitable for application to or for introduction into relevant orifice(s) of the body, e.g. the rectal, urethral, vaginal, aural, nasal or oral orifices. The composition may simply be applied directly on the part to be treated such as, e.g., on the mucosa, or by any convenient route of administration.\nCompositions which have proved to be of importance in connection with topical application are those which have thixotropic properties, i.e. the viscosity of the composition is affected e.g. by shaking or stirring so that the viscosity of the composition at the time of administration can be reduced and when the composition has been applied, the viscosity increases so that the composition remains at the application site.\nCompositions for oral use or for application to mucosa or skin\nSuitable compositions for use according to the invention may also be presented in the form of suspensions, emulsions or dispersions. Such compositions contains the active enamel substance in admixture with a dispersing or wetting agent, suspending agent, and/or one or more preservatives and other pharmaceutically acceptable excipients. Such compositions may also be suitable for use in the delivery of the active enamel substance to e.g. an intact or damaged mucosa such as the oral, buccal, nasal, rectal, or vaginal mucosa, or for administration to intact or damaged skin, or wounds.\nSuitable dispersing or wetting agents are, for example, naturally occurring phosphatides, e.g., lecithin, or soybean lecithin; condensation products of ethylene oxide with e.g. a fatty acid, a long chain aliphatic alcohol, or a partial ester derived from fatty acids and a hexitol or a hexitol anhydride, for example polyoxyethylene stearate, polyoxyethylene sorbitol monooleate, polyoxyethylene sorbitan monooleate, etc.\nSuitable suspending agents are, e.g., naturally occurring gums such as, e.g., gum acacia, xanthan gum, or gum tragacanth; celluloses such as, e.g., sodium carboxymethylcellulose, microcrystalline cellulose (e.g. Avicel(copyright) RC 591, methyl-cellulose); alginates and chitosans such as, e.g., sodium alginate, etc.\nSuitable examples of preservatives for use in compositions according to the invention are the same as those mentioned above.\nCompositions for use according to the invention may also be administered by the oral route. Suitable oral compositions may be in the form of a particulate formulation or in the form of a solid, semi-solid or fluid dosage form.\nCompositions for oral use include solid dosage forms such as, e.g., powders, granules, granulates, sachets, tablets, capsules, effervescent tablets, chewable tablets, lozenges, immediate release tablets, and modified release tablets as well as fluid or liquid formulations such as, e.g. solutions, suspensions, emulsions, dispersions, and mixtures. Furthermore, composition may be in the form of powders, dispersible powders, or granules suitable for preparation of an aqueous suspension by addition of a liquid medium such as, e.g. an aqueous medium,\nWith respect to solid dosage forms for oral (or topical use) a composition for use according to the invention normally contains the active enamel substance and any further active substance optionally in admixture with one or more pharmaceutically acceptable excipients. These excipients may be, for example,\ninert diluents or fillers, such as sucrose, sorbitol, sugar, mannitol, microcrystalline cellulose, starches including potato starch, calcium carbonate, sodium chloride, lactose, calcium phosphate, calcium sulfate, or sodium phosphate;\ngranulating and disintegrating agents, for example, cellulose derivatives including microcrystalline cellulose, starches including potato starch, croscarmellose sodium, alginates, or alginic acid and chitosans;\nbinding agents, for example, sucrose, glucose, sorbitol, acacia, alginic acid, sodium alginate, gelatin, starch, pregelatinized starch, microcrystalline cellulose, magnesium aluminum silicate, sodium carboxymethylcellulose, methylcellulose, hydroxypropyl methylcellulose, ethylcellulose, polyvinylpyrrolidone, polyvinylacetate, or polyethylene glycol; and chitosans;\nlubricating agents including glidants and antiadhesives, for example, magnesium stearate, zinc stearate, stearic acid, silicas, hydrogenated vegetable oils, or talc.\nOther pharmaceutically acceptable excipients can be colorants, flavouring agents, plasticizers, humectants, buffering agents, etc.\nIn those cases where the pharmaceutical composition is in the form of a solid dosage form in unit dosage form (e.g. a tablet or a capsule), the unit dosage form may be provided with a coating like one or more of the coatings mentioned below.\nIn those cases where the composition is in the form of a tablet, capsule or a multiple unit composition, the composition or the individual units or a tablet or a capsule containing the individual units may be coated e.g. with a sugar coating, a film coating (e.g. based on hydroxypropyl methylcellulose, methylcellulose, methyl hydroxyethylcellulose, hydroxypropylcellulose, carboxymethylcellulose, acrylate copolymers (Eudragit), polyethylene glycols and/or polyvinylpyrrolidone) or an enteric coating (e.g. based on methacrylic acid copolymer (Eudragit), cellulose acetate phthalate, hydroxypropyl methylcellulose phthalate, hydroxypropyl methylcellulose acetate succinate, polyvinyl acetate phthalate, shellac and/or ethylcellulose). Furthermore, a time delay material such as, e.g., glyceryl monostearate or glyceryl distearate may be employed.\nRectal and/or vaginal compositions\nFor application to the rectal or vaginal mucosa, suitable compositions according to the invention include suppositories (emulsion or suspension type), enemas, and rectal gelatin capsules (solutions or suspensions). Appropriate pharmaceutically acceptable suppository bases include cocoa butter, esterified fatty acids, glycerinated gelatin, and various water-soluble or dispersible bases like polyethylene glycols and polyoxyethylene sorbitan fatty acid esters. Various additives like, e.g., enhancers or surfactants may be incorporated.\nNasal compositions\nFor application to the nasal mucosa (as well as to the oral mucosa), sprays and aerosols for inhalation are suitable compositions according to the invention. In a typical nasal composition, the active enamel substance is present in the form of a particulate formulation optionally dispersed in a suitable vehicle. The pharmaceutically acceptable vehicles and excipients and optionally other pharmaceutically acceptable materials present in the composition such as diluents, enhancers, flavouring agents, preservatives, etc. are all selected in accordance with conventional pharmaceutical practice in a manner understood by the persons skilled in the art of formulating pharmaceuticals.\nDosages of enamel matrix, enamel matrix derivatives and enamel matrix proteins\nIn a pharmaceutical composition for use according to the invention on skin or mucosa, an active enamel substance is generally present in a concentration ranging from about 0.01% to about 99.9% w/w. The amount of composition applied will normally result in an amount of total protein per cm2 wound/skin/tissue area corresponding to from about 0.01 mg/cm2 to about 20 mg/cm2 such as from about 0.1 mg/cm2 to about 15 mg/cm2.\nThe amount applied of the composition depends on the concentration of the active enamel substance in the composition and of the release rate of the active enamel substance from the composition, but is generally in a range corresponding to at the most about 15-20 mg/cm2.\nIn those cases where the active enamel substance is administered in the form of a fluid composition, the concentration of the active enamel substance in the composition is in a range corresponding to from about 0.1 to about 50 mg/ml. Higher concentrations are in some cases desirable and can also be obtained such as a concentration of at least about 100 mg/ml.\nWhen the composition is applied to the oral cavity, the following doses are relevant:\nExperimental defect areas (in monkeys) in the oral cavity typically have a size of about 4xc3x972xc3x975-6 mm corresponding to 50 xcexcl or from about 0.025 to about 0.15 mg total protein/mm2 or about 2.5-15 mg/cm2. Usually up to 0,5 such as, e.g., 0.4, 0.3, 0.2 or 0.1 ml of a composition having a concentration of about 1-40 mg/ml such as, e.g., 5-30 mg/ml is applied.\nDefect areas in humans in the oral cavity and due to periodontal diseases typically have a size of about 5-10xc3x972-4xc3x975-10 mm corresponding to about 200 xcexcl and normally at the most about 0.5-1 ml such as about 0.2-0.3 ml per tooth is applied of a composition having a concentration of about 1-40 mg total protein/ml such as, e.g., 5-30 mg/ml is applied. 0.2-0.3 mg/ml corresponds to about 6 mg protein per 25-100 mm2 or about 0.1 mg/mm2 if calculated only on root surface. Normally an excessive volume is applied to allow coverage of all surfaces. Even a multilayer would only require a small fraction of the above-mentioned amounts.\nGenerally, about 0.1-0.5 ml such as, e.g., about 0.15-0.3 ml or about 0.25-0.35 ml of a composition comprising the active enamel substance is applied in defect volumes in extraction alveoli (holes after extraction of teeth). The concentration of the active enamel substance in the composition is normally about 1-40 mg total protein/ml such as, e.g., 5-30 mg/ml. When 0.3-0.4 ml is applied of such a composition for wisdom teeth, this volume corresponds to about 0.1 mg/cm2 (alveolus calculated as cylinder which radius 5 mm and height 20 mm).\nThe concentration of the active enamel substance in a pharmaceutical composition depends on the specific enamel substance, its potency, the severity of the disease to be prevented or treated, and the age and condition of the patient. Methods applicable to selecting relevant concentrations of the active enamel substance in the pharmaceutical composition are well known to a person skilled in the art and may be performed according to established guidelines for good clinical practice (GCP) or investigational New Drug Exemption (xe2x80x9cINDxe2x80x9d) regulations as described in e.g. International Standard ISO/DIS 14155 Clinical investigation of medical devices, 1994 and ICH (International Committee for Harmonisation): Harmonised tripartite guideline for good clinical practice, Brookwood Medical Publications, Ltd, Surrey, UK, 1996. A person skilled in the art would, by use of the methods described in standard textbooks, guidelines and regulations as described above as well as common general knowledge within the field, be able to select the exact dosage regimen to be implemented for any active enamel substance and/or selected other active substances and dosage form using merely routine experimentation procedures.\nIn other aspects the invention relates to methods for i) preventing and/or treating wounds, ii) decreasing infection and iii) preventing and/or treating inflammation, the methods comprising administration to a mammal in need of such a treatment an effective amount of an active enamel substance.\nAs will be understood, details and particulars concerning the use of an active enamel substance for the prevention and/or treatment of wound will be the same as or analogous to the details and particulars concerning the other use aspects (anti-bacterial and anti-inflammatory aspects) and the method aspects discussed above, and this means that wherever appropriate, the statements above concerning an active enamel substance, a preparation containing an active enamel substance, a pharmaceutical composition containing an active enamel substance, preparation of i) an active enamel substance, ii) a preparation containing an active enamel substance, ii) a pharmaceutical composition containing an active enamel substance, as well as improved properties and uses apply mutatis mutandis to all aspects of the invention."}
{"text": "Accumulation of solid and fibrous debris in drain pipes can cause backup of unsightly and unhygienic fluids in showers, kitchen sinks and lavatories. Clogged residential and commercial pipes are often cleared with abrasive chemical solutions which may harm individuals and the environment. The prior art discloses numerous systems employing traps, filters and strainers for preventing materials from clogging drain pipes.\nU.S. Pat. No. 6,088,843 discloses a drain strainer for installation into the vertical drain pipe of a bathing water drain for a shower or tub. The drain strainer includes a sieve receptacle having a substantially cylindrical basket, a bristle extension with a stem having a plurality of bristles, a support structure to suspend the sieve receptacle and bristle extension within the drain pipe.\nU.S. Pat. No. 6,487,729 B2 discloses a hair and fibrous shower trap that comprises a material trapping system loop structure. The loop structure is comprised of multiple, concentric circles fixedly attached to each other. The loop structure has a large opening at one end and a smaller opening at the other. The loop structure is dimensioned to fit pop-up, plunger, or push-pull drains for showers, bathtubs, or sinks. In one embodiment, the loop structure is called a salon trap which is shaped in such a manner as to fit within the screw-on opening to a sink trap.\nU.S. Pat. No. 4,490,862 discloses a drain filter for preventing hair and debris from entering a drain of a bathtub, shower stall, lavatory, sink and the like having a liftable stopper and a flange. The filter has a body shaped and dimension to be adapted to extend entirely around the lifted stopper periphery and extends around the external periphery of the drain. The body is fabricated of a filamentary material with openings to allow for the free passage of water down into the drain, and has a plurality of outward projections above the opening to entangle hair and debris.\nThese approaches do not present an optimal system for trapping and removing hair from a drain pipe. Solutions which employ a basket-like structure to trap fibrous materials may be circumvented if the basket becomes clogged with small debris. Debris that accumulates in the lowest point of the basket may impede flow down the drain and inhibit the trapping action of the basket. Furthermore, unpleasant odors may be emitted from the drain pipe if debris, such as hair, is trapped in the low point of basket placed at an insufficient depth in the drain pipe. Solutions that place a filter in only the trap region of a drain will likely result in blockage and ineffective capture of fibrous materials since the trapping action takes place at the lowest point instead of the vertical drain pipe. Furthermore, removal of such a filter may be cumbersome since the filter would have to be removed from the bottom of the trap region. Solutions that employ a one system approach do not provide redundancy in the event that the capturing mechanism is compromised. Solutions that require modifications to be made such as altering the pop-up stopper can be both costly, difficult and unsanitary to remove.\nIn light of the shortcomings in the prior art, a system for trapping is desired that employs an ensnaring systems that has a multiple backup ensnaring mechanisms in the event that the primary ensnaring mechanism is circumvented. The ideal system will be easy to remove from a drain pipe, disposable, and extend from the vertical drain pipe up to or beyond the trap, thereby providing a large surface area for ensnaring materials."}
{"text": "1. Field of the Invention\nThe present invention relates to resonators having ready-to-use conductor lines mounted with other circuit components on printed circuit boards, and more particularly relates to a resonator in which inductance or impedance may be varied uniformly in the longitudinal direction of a conductor line.\n2. Description of the Related Art\nA resonator, for example, for determining the oscillating frequency of a voltage-controlled oscillator, has been constructed with a strip conductor line directly formed on a printed circuit board on which other circuit components for the circuit construction are mounted.\nFIG. 7 illustrates a conventional resonator. In FIG. 7, a printed circuit board 21 has circuit components (not shown) constituting the voltage-controlled oscillator or the like, mounted thereon. A strip conductor line 22 is directly formed on this printed circuit board 21 and is connected to circuit components.\nFor example, when the oscillating frequency is adjusted, as shown in FIG. 8, the inductance or the impedance of the resonator is varied by locally changing the width of a slit 23 provided in the conductor line 22.\nAs described above, the conventional resonator, in which the conductor line 22 is used, is directly formed on the printed circuit board on which circuit components are mounted. Accordingly, in a case where conductor lines 22 having different widths or different lengths are required, since the design for each printed circuit board is different, a universal printed circuit board cannot be designed.\nSince the oscillating frequency is adjusted by locally changing the width of the slit 23, the inductance or the impedance of the conductor line 22 varies in accordance with the position of the conductor line 22 in the longitudinal direction thereof. In order to secure the required range of frequency change, the optimal position of the slit 23 must be selected. Moreover, the impedance locally increases in the periphery of the slit 23. Additionally, as the width of the conductor line 22 decreases, loss therein increases and the Q-factor therein decreases."}
{"text": "A silicon nitride film is used as an insulating film, a protection film or the like in various semiconductor devices. The silicon nitride film may be formed by a plasma CVD method using a silicon-containing gas such as silane (SiH4) and a nitrogen-containing gas such as nitrogen or ammonia as source gases (see, e.g., Japanese Laid-open Publication No. 2000-260767).\nIn the silicon nitride film formed by a conventional plasma CVD method, it is important to suppress stresses, i.e., tensile and compressive stresses, of the film, which have a bad influence on device characteristics. For instance, if a compressive stress of the silicon nitride film is large, stress migration may occur to cause disconnection of metal lines right under the film. Accordingly, the compressive stress needs to be reduced to prevent such disconnection. In the plasma CVD method, the direction of stress (tensile stress or comprehensive stress) of the silicon nitride film or the magnitude of stress depends on process conditions including pressure, temperature, gas species, etc. Thus, conventionally, process conditions are set such that a strong stress is not applied to the silicon nitride film, and a silicon nitride film having no stress is formed by a plasma CVD method (e.g., Maeta Gazuo ┌VLSI and CVD┘ Tenshoten, published on Jul. 31, 1997).\nRecently, attempts have been made to take advantage of the stress of silicon nitride film to improve the device characteristics. However, for example, a parallel plate type or inductively coupled plasma CVD apparatus utilizes plasma having a relatively high electron temperature. Accordingly, when film formation conditions such as high-frequency output, pressure and temperature are changed to introduce a high stress, plasma damage may occur in the silicon nitride film. Thus, it is difficult to obtain a high-quality silicon nitride film and it is also difficult to obtain a silicon nitride film having a high stress. Further, since there is a limitation in the plasma process conditions, it is also difficult to precisely control a stress."}
{"text": "A mobile agent can be written as a program that executes on a set of network hosts. The agent visits the network hosts to execute parts of its program. The agent may need, for example, to access information located on a given network host or there may be some preference to execute parts of its program on various network hosts.\nIn prior art systems, the sequence of hosts that the agent visits is statically pre-configured when an agent program is written. Furthermore, the agent in these systems must execute an explicit instruction to move to another remote network host. For example, U.S. Pat. No. 5,603,031, issued Feb. 11, 1997, entitled \"System and Method for Distributed Computation Based Upon the Movement, Execution, and Interaction of Processes in a Network,\" by White et al., describes a method for statically pre-configuring an agent's itinerary in a destination list composed of destination objects. Each destination object has a telename and a teleaddress preassigned to specific regions of the network. In the system disclosed in White et al., an agent must execute a \"go\" statement to move to another network host to access resources located at that host.\nThere are many other examples of this explicit agent relocation requirement. In Lange et. al., IBM Aglets Workbench--Programming Mobile Agents in Java, Proceedings of 1997 World Wide Computing and Its Applications, Japan, pp. 253-266, the authors suggest that an agent execute the \"dispatch\" statement to move to another network host. In Cardelli, A Language with Distributed Scope, Computing Systems, Vol. 8, No. 1, Winter 1995, pp. 27-59, Cardelli describes a system wherein object migration is achieved by explicitly copying the object (agent) state. Finally, in Chess et al., Itinerant Agents for Mobile Computing, available as IBM Research Report RC-20010, the authors state the need for explicit primitives and mobility support to make an agent move to another network host.\nSystems that require this explicit agent relocation suffer from increased programming complexity because the programmer must be aware of the location where a particular piece of code would execute. As a result, the application code has to be organized into location-sensitive blocks with explicit instructions for agent relocation demarcating the blocks. Location awareness distracts the programmer from focusing solely on developing application logic. This reduces programmer productivity.\nThus, there is a need for a system which permits a program (agent) to execute throughout a network without the need for explicit agent relocation, thereby improving programmer productivity."}
{"text": "1. Field of the Invention\nThe present invention relates to a silicon-containing triarylamine compound, a method of its production and a high-sensitivity high-durability photoreceptor for electrophotography incorporating said silicon-containing triarylamine compound as the charge transport material.\n2. Description of the Prior Art\nIn recent years, there have been noticeable improvements in copying machines and printers based on electrophotography. A wide variety of new models of different shapes, types and functions have been developed according to uses, and various photoreceptors for such new models have been developed accordingly.\nTraditionally, photoreceptors for electrophotography have been based mainly on an inorganic compound from the viewpoint of its sensitivity and durability. Examples of such inorganic compounds include zinc oxide, cadmium sulfide and selenium. However, these substances are often toxic, thus posing problems of disposal and environmental pollution. In addition, when using selenium, which offers good sensitivity, it is necessary to form a thin film on a conductive substrate by vacuum deposition, etc., which results in poor productivity and increased production cost.\nAmorphous silicon has recently drawn much attention as an inorganic photoreceptor free of environmental pollution, and research and development activities are in progress. However, the use of the plasma CVD method for film formation results in extremely poor productivity and increased photoreceptor production cost and running cost, though excellent sensitivity is obtained.\nOn the other hand, organic photoreceptors have the advantage of being free of environmental pollution since they are burnable, many of which permit easy mass production since they are coatable for thin films. These aspects are advantageous in that significant cost reduction is possible and a wide range of shapes can be obtained according to uses. It should be noted, however, that organic photoreceptors involve some problems to be solved as to sensitivity and durability; there are strong demands for the development of an organic photoreceptor with high sensitivity and high durability.\nVarious methods have been proposed to improve the sensitivity of organic photoreceptor, but the current main stream uses photoreceptors having a double-layered photoconductive structure in which functions are allotted to a charge generation layer and a charge transport layer. For example, the charge generated in the charge generation layer in response to exposure is injected to the charge transport layer, through which it is transported to the surface of the photoreceptor, where the surface charge is neutralized and an electrostatic latent image is formed thereon. Thus, in the double-layered photoconductive structure type in comparison with the single-layered structure type, the generated charge is less likely to be trapped, and the charge can be more efficiently transported to the surface of the photoreceptor (U.S. Pat. No. 2,803,541).\nAs the organic charge generation material for the charge generation layer, a compound which absorbs the energy of the irradiated light and efficiently generates a charge is selected and used. Examples of such compounds include an azo pigment (Japanese Patent Publication Open to Public Inspection No. 14967/1979), a metal-free phthalocyanine pigment (Japanese Patent Publication Open to Public Inspection No. 143346/1985), a metallophthalocyanine pigment (Japanese Patent Publication Open to Public Inspection No. 16538/1975) and a squarylium salt (Japanese Patent Publication Open to Public Inspection No. 27033/1978).\nAs the charge transport material for the charge transport layer, it is necessary to select a compound which offers a high efficiency of charge injection from the charge generation layer and a high charge mobility in the charge transport layer. To meet these requirements, a compound having a low ionization potential or a compound which easily releases a cationic radical is selected. Examples of such compounds proposed include a triarylamine derivative (Japanese Patent Publication Open to Public Inspection No. 47260/1978), a hydrazone derivative (Japanese Patent Publication Open to Public Inspection No. 101844/1982), an oxadiazole derivative (Japanese Patent Examined Publication No. 5466/1959), a pyrazoline derivative (Japanese Patent Examined Publication No. 4188/1977), a stilbene derivative (Japanese Patent Publication Open to Public Inspection No. 198043/1983), a triphenylmethane derivative (Japanese Patent Examined Publication No. 555/1970), and a 1,3-butadiene derivative (Japanese Patent Publication Open to Public Inspection No. 287257/1987).\nAs stated above, however, these organic photoreceptors are lower than inorganic photoreceptors in charge mobility, and their sensitivity remains unsatisfactory.\nAlso, in a series of electrophotographic processes of charging, exposure, development, toner transfer and discharging, the photoreceptor is subject to extremely severe conditions and its ozone resistance and abrasion resistance are of great concern. Although improvements in these durability properties are demanded, no satisfactory photoreceptors have been developed.\nAbove all, triarylamine has been investigated for a large number of derivatives, including monotriarylamine (Japanese Patent Publication Open to Public Inspection Nos. 195254/1982, 118147/1989 and 118142/1989), bis(diarylamino)benzene (Japanese Patent Publication Open to Public Inspection Nos. 144250/1980, 118144/1989 and 118146/1989) and bistriarylamine (Japanese Patent Publication Open to Public Inspection Nos. 27033/1978, 35140/1981, 52756/1981 and 142647/1989).\nHowever, monotriarylamine has drawbacks including insufficient charge mobility and a lack of stability, though it is generally easy to synthesize. As for bis(diarylamino)benzene and bistriarylamine, they have drawbacks such as a lack of stability and difficult synthesis, though they offer a fair level of charge mobility. Also, these triarylamine derivatives all have been suggested to pose a problem regarding a lack of compatibility with binder resin.\nAs stated above, conventional organic photoreceptors such as triarylamine derivatives have many drawbacks, and there are strong demands for improvements to overcome these drawbacks in the relevant technical field.\nIt should be noted that a silicon-containing triarylamine compound which is similar to the present invention is disclosed in U.S. Pat. No. 2,960,517, but it is a monosilyltriarylamine compound and is unsuitable for use as a charge transport material for photoreceptor for electrophotography because it is liquid at normal temperature. In addition, it is produced by reaction of a bromophenylalkylsilane compound and an alkali metal salt of amine, and these starting materials are not easily available and this reaction requires severe conditions both for reaction temperature and for reaction time."}
{"text": "The present invention relates generally to time management devices for use by pilots and, more specifically, to a portable computer system that allows pilots to more efficiently and effectively manage the time-oriented and other aspects of their flights.\nThere are basically two areas in which pilots must adhere to time-sensitive procedures. The first are precision flight maneuvers in which time is critical to remain within certain Federal Aviation Administration (FAA) limits, and the other is to monitor flight progress and aircraft subsystems on a periodic basis to ensure safe and proper operation of the aircraft. A typical flight may be divided into the following phases: pre-flight, enroute, approach, holding and post-flight. The enroute route phase is made up of a series of legs between various waypoints. Pilots must monitor the actual time taken to fly each leg to obtain accurate fuel consumption information. During the approach and holding phases of the flight, the pilot typically must perform a series of precise, time-dependent maneuvers.\nBefore taking off, many pilots prepare a handwritten flight plan. This flight plan typically includes the waypoints that are to be flown during the flight. The handwritten flight plan is often attached to a knee board that is secured to the pilot\"\"s leg for reference during the flight. To the extent it spans multiple pages, handwritten flight plans can be difficult to consult during the flight. They can also be difficult to read.\nFurthermore, during each phase of the flight, pilots are often confronted with multiple demands on their attention and concentration. For example, in addition to flying the aircraft, pilots must typically handle all navigation and communication duties, monitor weather conditions, monitor the aircraft\"\"s fuel supply and perform other such tasks. Pilots also must continually review and check cockpit indicators and gauges to monitor the condition and operation of the aircraft\"\"s many subsystems. In particular, most aircraft have one or more cockpit displays that reveal the operating condition of the engine, the hydraulic systems, the electrical systems, the fuel systems, the landing gear systems, the auxiliary power units, if any, etc. Pilots flying in inclement weather, in close proximity to other aircraft or in other demanding conditions often neglect to continually scan their indicators and gauges. As a result, pilots may overlook the early warnings of an impending failure until it is too late to take corrective action.\nWhile enroute, the pilot may also encounter an emergency situation. Most aircraft have emergency checklists identifying the corrective action to be taken in response to many different types of emergencies. Although these procedures are often contained in a loose-leaf, print format that is designed for relatively easy access, it can be difficult and/or time consuming to locate the specific pages corresponding to the particular emergency being faced. Flipping between multiple pages of the emergency checklist is similarly awkward especially where the pilot is busy controlling the flight characteristics of the aircraft as a result of the emergency.\nIn addition, the pilot may be interrupted during execution of the emergency checklist. For example, a call may come over the communications system that must be responded to or some other action may need to be taken. When the pilot returns to the checklist, he or she may forget the point at which he or she was interrupted. That is, the pilot may not remember which portions of the checklist have been completed and which portions still need to be completed.\nDespite the importance of time-management during flight, the only devices currently available are conventional stopwatch timers and sweep second-hand stopwatches and clocks. Accordingly, a need exists for a more comprehensive and user-friendly device to assist pilots in time-management and other flight-related tasks.\nBriefly, the present invention is directed to a system and method that facilitates the management of time-oriented and other flight-related tasks. In the illustrative embodiment, the system includes a hand-held computer having a microprocessor, a memory and a display/input screen. The microprocessor, which includes or is coupled to a timer, preferably executes one or more novel application programs that are configured to request information from and display or provide information to the pilot at appropriate times during a flight. The application program is preferably menu-driven so that the pilot may easily navigate among the various displays and retrieve desired information or initiate desired functionality, such as activating one-or-more-count-Tupor count-down timers or displaying one or more-electronic-checklists The system may also include a visual, tactile and/or aural warning element.\nIn a first aspect, the system, including the application program, is configured to display one or more bulletins at pre-arranged times during flight. More specifically, during the preflight phase, the pilot preferably programs the system to generate and display one or more bulletins prompting the pilot to check the status of specific aircraft subsystems or to perform some other flight-related tasks. The bulletins can be programmed to occur at a single preset time during the flight or they can be set to occur periodically throughout the flight. For example, the system can be programmed to display a first bulletin directing the pilot to check the engine oil and fuel systems every seven minutes, while a second bulletin prompts the pilot to check the generator or electrical subsystem only every fourteen minutes. As the flight progresses, the system automatically interrupts the current application or process and displays the bulletins at the programmed time intervals. Once the displayed task has been performed, the pilot preferably acknowledges its completion by pressing an xe2x80x9cacknowledgementxe2x80x9d key or button which may be displayed on the screen and/or may be remotely located from the system, e.g., a thumb switch coupled to the system. Following the acknowledgement, the application program preferably causes the next bulletin to be displayed at the appropriate time. If no acknowledgment is received after a pre-set time, the application program may activate the warning element.\nIn a second aspect, the system facilitates the creation of pilot reports (PIREPS). that can be transmitted by the pilot to an air traffic controller (ATC). PIREPS preferably follow a specific, FAA-approved format (e.g., sequence of information) and utilize a number of abbreviations that can be difficult to remember. Because of the complexity of the format and abbreviations, many pilots do not provide PIREPS. The application program is preferably configured to display one or more windows or menus, upon request, that solicit particular information from the pilot, including location, weather, etc., in an easy-to-read and easy-to-understand format. The program includes a PIREP conversion engine that translates the information provided by the pilot into a PIREP-compatible format for display. The pilot can then simply read the system generated and displayed PIREP over the aircraft\"\"s communication subsystem to the ATC.\nIn a third aspect, the application program provides one or more novel timer displays to facilitate flying precision maneuvers, such as holding and approach maneuvers. Specifically, the application program is configured to generate an approach timer window that includes a settable pending timer field or window and an active timer field or window. Within the pending timer field, the pilot preferably enters the time value associated with flying the next leg or segment of the approach. When the pilot starts this leg, e.g., crosses the initial fix point, the pilot selects a start button. In response, the application program copies into the active timer field the time value that was entered in the pending timer field and begins a count-up or a count-down based on that time value. During execution of the first leg, the pilot can enter the time associated with the next approach leg in the pending timer field. At the end of the first leg, the pilot can again enter the start button causing the program to copy the newly entered time value from the pending timer window into the active timer window and to begin the count up or the count down. This process can be repeated until the pilot lands the aircraft."}
{"text": "Technical Field\nThe present invention relates to an electric power supply system.\nBackground Art\nHeretofore, there has been proposed an electric power supply system in which a direct current power supply such as a battery and fuel cell mounted on an electric vehicle such as an electric automobile and a fuel cell automobile is used to supply electricity to household electrical devices (for example, refer to Japanese Unexamined Patent Application, First Publication No. 2006-325392).\nThe electric power supply system disclosed in Japanese Unexamined Patent Application, First Publication No. 2006-325392 comprises: a vehicle having a means for supplying electric power to the outside of the vehicle; a stationary fuel cell system provided with an inverter; a load device that receives electric power supply from the stationary fuel cell system; and a system power supply that supplies electric power to the stationary fuel cell system. In the event of a power outage, this electric power supply system connects the vehicle and the stationary fuel cell system, and supplies electric power from the vehicle to the load device via the inverter of the stationary fuel cell system.\nAs a fuel cell, there is known a fuel cell stack in which a membrane electrode assembly is formed by arranging an anode electrode and a cathode electrode on either side of a solid polymer electrolyte membrane (hereunder, referred to as “electrolyte membrane”), arranging a pair of separators on either side of this membrane electrode assembly to form a flat unit fuel cell (hereunder, referred to as “unit cell”), and then stacking a plurality of these unit cells together to form a fuel cell stack. In this fuel cell, hydrogen ions produced by a catalytic reaction at the anode pass through the electrolyte membrane and move toward the cathode. There, they react with the oxygen in the air, giving rise to an electrochemical reaction and the generation of electric power.\nThe fuel cell described above generates heat as electric power generation is performed, and therefore, the generated water produced as a result of the electric power generation in the fuel cell is likely to vaporize. The generated water that has vaporized (water vapor) is discharged together with cathode off-gas and anode off-gas, and as a result, the electrolyte membrane of the fuel cell becomes dry. If the fuel cell becomes excessively dry (hereunder, referred to as “dry-up condition”), there is a problem in that the power generation performance of the fuel cell becomes reduced, and this consequently leads to deterioration in the electrolyte membrane.\nTherefore, the fuel cell system is provided with a cooling device for cooling the fuel cell which generates heat as power generation is performed. The cooling device is formed with a coolant that circulates in the fuel cell and absorbs heat, a radiator for releasing heat from the coolant, and a radiator fan that blows air to the radiator.\nIncidentally, in general, cooling devices of fuel cells and control thereof are designed in consideration of a vehicle in a state of traveling.\nSpecifically, they are designed such that in those cases where the traveling speed of the vehicle is high and the amount of electric power being generated by the fuel cell is high, traveling air stream is introduced into the radiator and the radiator fan is rotated at a high rotation speed, to release the heat of the coolant flowing through the radiator. Moreover, they are designed such that in those cases where the traveling speed of the vehicle is low and the amount of electric power being generated by the fuel cell is low, traveling air stream is introduced into the radiator and the radiator fan is rotated at a low rotation speed, to release the heat of the coolant flowing through the radiator. As a result of this, the fuel cell is appropriately cooled to suit the traveling state of the vehicle (to suit the electric power generation state of the fuel cell), and it is therefore possible to prevent the electrolyte membrane from becoming dry when the vehicle is traveling.\nFurthermore, the fuel cell system is designed so that rotation of the radiator fan stops when the vehicle is stopped and the electric power generation of the fuel cell is stopped. As a result, wasteful electric power consumption by the radiator fan can be prevented."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, and more particularly relates to a semiconductor device that includes a DLL (Delay Locked Loop) circuit having a coarse adjusting unit and a fine adjusting unit.\n2. Description of Related Art\nSome of high-speed semiconductor memory devices such as a DRAM (Dynamic Random Access Memory) use a DLL circuit to accurately match an output timing of read data with a clock signal. The DLL circuit generates an internal clock signal accurately phase-controlled and, in some cases, employs a coarse adjusting unit that controls a phase of the internal clock signal at a relatively large adjustment pitch and a fine adjusting unit that controls the phase of the internal clock signal at a relatively small adjustment pitch (see Japanese Patent Application Laid-open Nos. 2009-284266 and 2009-21706).\nHowever, when an adjustment code supplied to the fine adjusting unit is changed, jitters may be caused in the internal clock signal. It is therefore desired that the semiconductor device can decrease jitters caused by changes in the adjustment codes."}
{"text": "This invention relates generally to ferroelectric films. More particularly, the present invention relates to a ferroelectric film having improved electrical characteristics and the corresponding fabrication method for producing the ferroelectric film.\nFerroelectric films are typically used as the dielectric material for a ferroelectric capacitor in a ferroelectric memory cell. For proper memory operation, it is critical that the ferroelectric film achieves desirable electrical performance evidenced by the ability to liberate a detectable charge in response to the application of an externally applied electrical field. It has been shown in the laboratory that crystal growth characteristics and film orientation are critical to achieving this desirable electrical performance. Typically, the ferroelectric film is deposited on a conducting electrode surface, such as the bottom electrode of the ferroelectric capacitor. The bottom electrode can be fabricated from a variety of films such as: platinum, iridium, iridium oxide, ruthenium oxide, palladium, as well as other noble metals and their oxides, or other suitable conductive materials known in the art. Each of these electrodes has a unique surface with changing roughness, and/or conductivity, which dramatically impacts the sticking coefficients for the various components of a subsequently deposited ferroelectric film. As the sticking coefficients vary, the composition of the ferroelectric film varies. As the electrode surface begins to accept a layer of ferroelectric film, the sticking coefficients for the constituents change and dictate the composition of the bulk film. Therefore, the composition required for optimum electrical performance cannot be obtained using a single fixed set of sputtering conditions or solution chemistry.\nFor subsequent comparison to the improved ferroelectric film of the present invention, a typical prior art ferroelectric PZT film (lead zirconate titanate) having a thickness of about 2400Angstroms and deposited with a single step deposition has roughly 75% of the crystal domains oriented in the &lt;111&gt; crystal orientation. The \"switched charge\" liberated using a three-volt pulse [\"Qsw(3 v)\"] was measured at about 21 micro-Coulombs per centimeter squared. The applied voltage at which the switched charge of the ferroelectric capacitor is 90% saturated [\"V(90%)\"] was measured at about 4.5 volts.\nWhat is desired is a ferroelectric film having optimized electrical performance in which specific adjustments in deposition conditions are made to compensate for the electrode variables described above and to dictate a desired film composition profile."}
{"text": "1. Field of the Invention\nThe present invention relates to a planing boat, particularly of small type, equipped with an airbag system.\n2. Related Art\nRecently, it is general for a four-wheel car to be provided with an airbag system as means for relaxing an impact on riders in a collision.\nThe airbag systems for four-wheel cars are based on the premise of using with a seatbelt and used for relaxing an impact on the riders in a car collision, and the seatbelt is for absorbing a force that throws the riders forwards.\nAs marine sports become popular, a collision between small planing boats such as motor boats and wet bikes (or personal watercrafts) increases. Accordingly, airbag systems for relaxing the impact on the riders at a boat collision have been researched and developed to be mounted.\nAn example of the airbag system provided for motor boats is disclosed in Japanese Patent Laid-open (KOKAI) Publication No. HEI 7-108985, which is equipped with an airbag on the deck obliquely behind the cockpit surrounded by a cowling, this airbag being inflated manually when the rider feels or detects danger.\nAlthough there is no example of airbags designed specifically for a wet bike in which a cockpit is exposed to the outside, there is another similar example of a wearable airbag suitable for protecting motor-bike riders (for example, refer to Japanese Patent Laid-open Publication No. HEI 9-66789.\nHowever, a small planing boat such as wet bike sometimes utilizes weight shift of the rider on a saddle-type driver's seat to control the operation, and in such a planing boat, no seatbelt is equipped. Accordingly, the airbag system for the four-wheel car based on the premise of utilizing the seatbelt is not suitable for such small planing boats. In addition, because of no seatbelt equipped, at a time of collision, a large force directed to the outside of the boat is applied to the rider in the small planing boat, being dangerous."}
{"text": "Fibrous structures, such as sanitary tissue products, have utilized oils, especially silicone oils, in the past to provide surface softening. However, consumers of fibrous structures continue to desire even more softness that is not deliverable solely by surface softening techniques and/or silicone oils. In other words, there is still an existing need for fibrous structures that exhibit even more softness than what can be delivered by surface softening techniques and/or silicone oils."}
{"text": "This invention relates to optical fiber units for use in an optical fiber cable.\nAn optical fiber cable usually comprises a plurality of elongated optical fiber units wound spirally within an outer sheath. Each unit comprises an optical fiber telescoped within and spaced inwardly from an outer tubular jacket so that an annular space is defined between the fiber and the jacket.\nIn some applications, the cable is exposed to high-pressure hazardous gases which must be blocked against progressing through the interior of the cable. For example, United States Department of Energy specifications for certain optical fiber cables require that the cable be capable of blocking flow at end pressures as high as 125 p.s.i. In part, gas blockage is achieved by filling the annular space between the fiber and jacket of each optical fiber unit with a material which establishes a pressure-tight seal between the fiber and the jacket.\nPrevious materials which have been used as fillers include non-silicone materials such as polyurethane jelly, photo-curable polymers, semi-liquid organic polymers, and a material sold under the trade designator \"INDOPOL\". Difficulty has been encountered with those materials in that prior materials which are sufficiently viscous to block flow at the specified gas pressure excessively restrict free movement of the fiber in the jacket and thus cause undesirable changes in certain optical properties of the fiber when tension is applied to the cable. On the other hand, those prior materials with a relatively low viscosity enabling substantially free movement of the fiber are not capable of resisting high gas pressures when such materials are used to fill an annular space of sufficiently large cross-sectional area to impart desired properties to the fiber unit. Moreover, the viscosity of the previously used materials changes significantly over the wide temperature ranges to which the fiber unit may be exposed, such materials typically becoming excessively stiff at low temperatures and excessively fluid at high temperatures."}
{"text": "The present disclosure relates to a method for manufacturing a magnetic coil for a magnetic actuator, particularly for an injector. The disclosure further relates to a magnetic coil for a magnetic actuator and to an injector, particularly a fuel injector nozzle for a combustion engine, comprising such a magnetic actuator.\nA common fuel injector nozzle is for instance known from EP 0 132 623 A2. The injector nozzle is provided with a magnetic actuator comprising an induction coil which is encapsulated and arranged in a fashion fixed to a housing.\nInjectors for the fuel supply for combustion engines are generally known. Injectors of that kind are also referred to as injector nozzles. Modern injectors are provided with mechatronical actors, for instance with magnetic actuators, so as to control the opening and closing of the nozzles in a high-precision fashion. Modern injection systems for diesel engines are configured for maximum injection pressures of several hundred bars up to more than 1000 bar or even up to 2500 bar.\nInjectors comprising so-called piezo-actuators are known. Nevertheless, also injectors comprising magnetic actuators are widely used. Injectors of that kind are commonly provided with a magnetic coil which cooperates with a movable armature for controlling the injector.\nActuators for injector nozzles must be manufactured in a high-precision fashion so as to be able to provide a respectively desired injection volume in a repeatedly accurate manner. Accordingly, there are frequently very tight tolerance specifications for the production. Further, it must be noted that the utilized components of the injector have to be considerably fuel-resistant or propellant-resistant so as to be resistant to the fuel which flows through the injector. The fuel may for instance involve diesel fuel, gasoline fuel, ethanol, light oil, heavy oil, kerosene or mixtures thereof.\nThe frequency at which the actuators of the injectors are controlled typically corresponds to the speed of rotation of the combustion engine. However, there are also injectors known that enable a plurality of injections event in a single combustion cycle. This may involve a further increased load to the injector and to the actuator thereof. Further, injectors are regularly fixedly attached to a cylinder head or a similar housing component of the combustion engine. This involves that the injector and the components provided therein are exposed to heavy vibration exposures, temperature variations and similar environmental impacts. Accordingly, there is a general need to design injectors in a robust fashion so as to ensure a specified minimum lifespan (for instance defined by operating cycles and such like).\nFor injectors that utilize magnetic actuators, it is a common measure to at least sectionally overmold and/or pot a coil winding of the magnetic coil of the actuator. In this way, the sensitive coil winding may be protected against environmental influences. Generally, such a magnetic coil comprises a coil body made from plastic material which supports a copper winding. For instance, the coil carrier that is provided with the coil winding is arranged in a magnet pot which defines a yoke of the magnetic actuator. By sealing, potting or overmolding the coil carrier that is received at or in the magnet pot, the coil winding is at least sectionally overmolded and/or potted with a (temporarily) fluid filler material.\nThe armature of such a magnetic actuator which cooperates with the magnet pot and the coil winding is frequently arranged as a plate armature and/or in a similar fashion. A lifting distance of the armature is frequently in a range of a few hundredths of a millimeter up to about 0.1 mm. As a result, even deviations in the range of a few micrometers at the seating surface and/or the abutment surface of the armature have a great influence on the function of the injector. A design goal may therefore involve to avoid, as much as possible, deviations of the armature stroke that is defined by the design.\nIt has been observed that the required potting and/or the required overmolding of the coil winding may in this respect have a negative impact. Generally, the overmolding and/or potting of the coil winding is performed in connection with the magnet pot, that is, in a state when the coil carrier that is provided with the coil winding is already arranged at the magnet pot in the desired final position. This involves that a molding form and/or potting form for such a procedure has to cooperate with the magnet pot in a high-precision fashion so as to ensure a sufficient sealing. If this is not the case, excess molding material (flash material) may be formed which involves events wherein the filler material (molding compound or potting compound) enters regions which actually shall not be entered for design and/or functional reasons. Even though there exist established solutions for avoiding such excess molding material, these solutions always involve increased production efforts. For instance, the tool which is used for overmolding or potting may be manufactured at an increased accuracy, and the magnet pot which cooperates with the tool may be also provided with an increased precision in a portion thereof that is arranged for contacting such a tool. However, this involves an increased production effort and also an increased wear on the side of the tool.\nBy contrast, in cases where no complete sealing of the magnet pot by the injection molding or potting tool may be ensured, excess molding material (flash material) may be present. Hence, in these cases, an increased scrap/waste or functional deficiencies have to be expected. Excess molding material may be generally referred to as residual flash.\nIn view of this, it is a first object of the present disclosure to provide a method for manufacturing a magnetic coil for a magnetic actuator which results in an improved quality of the potting and/or molding.\nIt is a further object to provide a method for manufacturing a magnetic coil for a magnetic actuator which involves an improved reproducibility and process reliability.\nIt is a further object to provide a method for manufacturing a magnetic coil for a magnetic actuator which prevents excess flash molding material to a great extent or even entirely.\nMore particularly, it is a further object to provide a method for manufacturing a magnetic coil for a magnetic actuator, wherein residual flash of the filler material at functional surfaces, particularly at contact surfaces or seating surfaces for the armature may be avoided to a considerable extent.\nIt is a further object to provide a method for manufacturing a magnetic coil for a magnetic actuator that is arranged to utilize robust tools/molds for the potting or overmolding which may achieve a great service life and/or lifespan.\nIt is a further object to provide a method for manufacturing a magnetic coil for a magnetic actuator that enables a reduction of efforts for quality management and a reduction of scrap production.\nIt is a further object to provide a method for manufacturing a magnetic coil for a magnetic actuator that enables the production of injectors which provide an increased operation performance and which are arranged to be controlled in a highly-accurate fashion.\nIt is a further object to present exemplary arrangements of magnetic coils and of injectors that are provided with respective magnetic coils that may profit from exemplary embodiments and aspects of the manufacturing method."}
{"text": "Typically, an exterior rearview mirror assembly includes a mirror actuator that is operable to adjust a mirror reflective element relative to the mirror casing so as to adjust a rearward field of view of the driver of the vehicle. Powerfold exterior mirror assemblies include an actuator that pivots or folds the mirror casing relative to the side of the vehicle."}
{"text": "Adhesives for optical films require durability so that foaming, lifting, peeling, etc. of optical films should not occur even in a high-temperature high-humidity atmosphere when optical members are bonded to adherends, and require light leakage prevention property so that the adhesives may flexibly follow dimensional changes of optical films in a high-temperature atmosphere so as not to cause light leakage in the case where two polarizing films that are optical films are bonded with the adhesive in such a manner that their polarization axes are at right angles to each other and allowed to stand at high temperature or at high temperature and high humidity. Further, rework property is also required in order that an optical member (member having an adhesive on an optical film) may be re-released from an adherend such as a liquid crystal panel without contaminating the adherend even in the case where the optical film is bonded to the adherend at the deviated position.\nAs such adhesives for optical films, acrylic-based adhesives have been mainly used in the past, and as an adhesive satisfying the above requirements, an acrylic-based adhesive obtained by blending an acrylic polymer of high-molecular weight with a medium- or low-molecular weight substance and crosslinking the resulting blend is known. In this acrylic-based adhesive, by crosslinking the acrylic polymer of high-molecular weight, cohesive force is enhanced to impart durability so as not to cause lifting and peeling, and by introducing the medium- or low-molecular weight substance, occurrence of light leakage is inhibited.\nHowever, with recent increase of sizes of displays, light leakage attributable to dimensional changes of optical films, particularly polarizing films, in a high-temperature atmosphere becomes a more serious problem, and the adhesives for optical films have been desired to have higher flexibility. In severe conditions, further, there occurs a problem that the low-molecular weight substance added for enhancing light leakage prevention property bleeds out to contaminate the adherend in the re-release process or to cause peeling.\nFor example, technique to increase a ratio (Mw/Mn) of the weight-average molecular weight (Mw) to the number-average molecular weight (Mn) to not less than 4 in order to improve heat durability and stress relaxation property has been disclosed (patent document 1). However, by merely increasing the Mw/Mn ratio to not less than 4, the stress relaxation property required for large displays is insufficient.\nIt is also known to enhance adhesion to plastics in a high-temperature high-humidity atmosphere by using a mixture of a (meth)acrylic-based polymer having a weight-average molecular weight of not less than 200,000 and a (meth)acrylic-based polymer having a weight-average molecular weight of less than 200,000 (patent document 2). In this method, however, the component having a lower weight-average molecular weight bleeds on the surface, so that problems of contamination in the re-release process and peeling under severe working conditions could not been solved.\nFurther, it is also known to user as an adhesive for optical films, a polymer in which an aromatic group-containing monomer is copolymerized (patent document 3 and patent document 4). However, its object is to control the refractive index of the adhesive in order to inhibit interfacial reflection between an optical member, such as a polarizing plate or a retardation plate, and the adhesive and between the adhesive and an adherend. Since the adhesive is designed so as to exhibit strong adhesion immediately after bonding, there is a problem that the adhesive is poor in a property that re-bonding is possible even if wrong bonding is made (rework property).\nMoreover, it is also known to use, as an adhesive for lowly polar films, a high-molecular weight polymer containing an alicyclic monomer or an aromatic group-containing monomer as a copolymer component (patent document 5). The adhesive, however, is designed so that the value (Mw/Mn) obtained by dividing the weight-average molecular weight of the copolymer by the number-average molecular weight thereof may become relatively small, and therefore, when the adhesive is used for a polarizing film having a large bond area or the like, there resides a problem that a stress due to dimensional change of the polarizing film or the like cannot be completely relaxed, and hence, sufficient light leakage prevention property is not obtained. Further, even if the alicyclic monomer is copolymerized, the light leakage prevention property is insufficient.\nPatent document 1: Japanese Patent Laid-Open Publication No. 341141/2002\nPatent document 2: Japanese Patent Laid-Open Publication No. 107507/2002\nPatent document 3: Japanese Patent Laid-Open Publication No. 173656/2002\nPatent document 4: Japanese Patent Laid-Open Publication No. 13029/2003\nPatent document 5: Japanese Patent Laid-Open Publication No. 053976/2005\nIt is an object of the present invention to provide an acrylic-based adhesive composition for optical films, which not only relaxes a stress due to dimensional change of a member having a particularly large bond area, such as a polarizing film, a retardation film or an elliptically polarizing film, and prevents contamination in the re-release process and peeling or foaming under high-humidity high-temperature conditions but also inhibits color nonuniformity due to light leakage, and to provide an optical member using the adhesive composition."}
{"text": "The present invention relates to portable enclosures, and in particular to a roll-formed frame arrangement therefor.\nPortable closures or shelters, such as those disclosed in U.S. Pat. Nos. 3,335,535 and 3,952,463 (which are incorporated by reference herein) have become increasingly popular in recent years. Such structures typically comprise a series of screened wall panels, which are interconnected along their upstanding edges by flexible hinges. The wall panels are arranged in a closed polygon configuration, and support a dome-shaped, fabric roof. A door is provided in the structure to access the interior, and solid or opague panels can be installed over the screens to keep out gentle rain, wind, insects, etc, as well as to provide privacy.\nThe frames for the wall panels are generally constructed from extruded aluminum, with channels disposed along the side edges in which an I-shaped flexible hinge is mounted. The extruded frames require a relatively thick sidewall, in the nature of 0.045 inches, thereby resulting in high material costs, as well as expensive transportation and storage. The flanges or beads along the side edges of the flexible hinge are rectangular, and are typically fitted rather tightly into the mating frame channels, so as to form a reasonably weathertight seal therebetween. Since the flexible hinges are rather long, the assembly of the hinges into the frame channels by pulling the hinge beads therealong is a difficult task which can easily result in damage to the flexible hinges.\nRetainers are used at the ends of each of the wall panels to hold the fastening cables in position. Heretofore, the legs of the retainer were manually compressed, and then inserted into the frame end. Wind and other forces acting on the roof and wall panels of the shelter are transmitted to the cables, and pull outwardly on the retainers. Prior to the present invention, separate left and right hand retainers were required to assemble the shelter, and the same had a tendency to inadvertently unlock from the frame when the shelter was exposed to gusts of wind, or other similar extreme forces."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device. More particularly, the present invention relates to a liquid crystal display (“LCD”) device capable of improving viewing angles.\n2. Description of the Related Art\nA LCD device can change the arrangement of liquid crystal molecules interposed between two substrates, in response to an electric field applied thereto, and display images. An image quality of a LCD device depends on viewing angles due to arrangement of liquid crystal molecules. A LCD device is able to display good images within a particular range of viewing angles which have a contrast ratio of more than 10:1. A display device for a desktop monitor, for example, has a viewing angle of greater than 90 degrees.\nThe contrast ratio is a measure of a display system, defined as the ratio of the luminance of the brightest color (e.g., white) to that of the darkest color (e.g., black) that the system is capable of producing. When a LCD device is able to display a darker and/or a brighter color than a color that a common system is capable of producing, the contrast ratio is increased. In order to display the darker color, a LCD device may decrease the light leakage of a LCD panel, adopt a normally black mode and/or decrease the reflectivity of a black matrix. Conversely, in order to display the brighter color, a LCD device may be designed to have a high aperture ratio for illuminating more light from a backlight.\nViewing angles according to azimuthal angles, having the specific contrast ratio, may be enhanced by a compensation film and/or multi-domains. A wide viewing angle mode, such as a multi-domain vertical alignment (“MVA”) mode and a patterned vertical alignment (“PVA”) mode, has been developed in order to improve the viewing angles according to azimuthal angles in aspect of multi-domains."}
{"text": "Real space transfer (RST) devices are known to the art. See, for instances, \"Heterojunction Band Discontinuities: Physics and Device Applications\", F. Capasso et al., editors, Elsevier 1987, especially pages 513-537, incorporated herein by reference.\nKnown RST devices are a transistor variously called charge-injection transistor (CHNIT) or negative resistance field effect transistor (NERFET), and the hot-electron erasable programmable random access memory (HE.sup.2 PRAM). See, for instances, U.S. Pat. No. 4,903,092, also incorporated herein by reference.\nBriefly, the transistor is a three-terminal device based on real-space transfer of hot electrons from a first to a second conducting region. The two conducting regions are separated by a barrier region and are contacted independently, with one of the conducting regions (referred to as the \"channel\") having two surface contacts (frequently referred to as \"source\" and \"drain\"). Application of a source-to-drain bias V.sub.sd leads to a heating of channel electrons and consequent charge injection into the second conducting layer. The channel thus acts as a hot electron emitter and the second conducting layer as a collector. This terminology will be used herein. The above discussed transistor shows a strong negative differential resistance in the source-drain characteristic (the NERFET action) and an efficient control of the injection current (I.sub.c) by the source-drain voltage (CHINT action).\nHE.sup.2 PRAM comprises, in addition to the above described elements, a \"deep\" drain that contacts both the emitter and the collector.\nA logic circuit that comprises prior art RST devices is also known. For instance, on page 520 of the above referenced monograph is disclosed a logic circuit comprising two NERFETs.\nThose skilled in the art are well aware of the desirability of having available devices that have novel operating characteristics, since such devices may make possible attainment of previously unachievable results, or may result in more economical attainment of some desired results. For instance, prior art implementation of basic logic functions (e.g., AND, NOR) generally requires a multiplicity of active elements (typically transistors). It would clearly be desirable, both from an economic and from a performance point of view, to be able to implement such logic functions with logic elements that comprise only a single active device, and thus reducing the device count and the number of gate delays. As a second example, it would be highly desirable to have available a device that can carry out a logic function which in the prior art required two or more separate logic elements. As a third example, it would be very desirable to have available a device that can carry out either a first or a second logic functions. This application discloses a novel device that has these and other advantageous attributes. It also discloses articles that utilize the novel device."}
{"text": "1. Field of the Invention\nThis application relates to an analog switch circuit, and more particularly, to the optimizing circuit characteristics according to the ON/OFF status of the analog switch circuit.\n2. Description of the Related Art\nSince a metal-oxide-semiconductor field-effect transistor (MOSFET) can easily be downsized and have superior element characteristics as well as readily-integrated process, a MOSFET is usually disposed as an analog switch when a communication system is in need, where the ON/OFF status of the analog switch can be controlled by a bias applied on a gate electrode of the MOSFET. When the analog switch is ON, a carrier channel (N-type or P-type) forms between a drain electrode and a source electrode of the MOSFET, which is equivalent 20 to a resistor present between the drain electrode and the source electrode. When the analog switch is OFF, no carrier channel or only narrow carrier channel exists between the drain electrode and the source electrode of the MOSFET, which is equivalent to a capacitor present between the drain electrode and the source electrode.\nIn addition to the gate, drain, and source electrodes, the MOSFET may include a body electrode for further tuning. In conventional designs, only one bias is applied on the body electrode, or only one resistor with constant resistance is connected to the body electrode. However, according to the description above, an ON switch and an OFF switch require two completely different element characteristics. Hence, only providing fixed bias to the body electrode or connecting the body electrode to the resistor with fixed resistance cannot further optimize the analog switch, and better quality for signal transmission cannot be obtained."}
{"text": "Quantum dots (QDs) provide many advantageous features that include high quantum yield, broad absorption spectra, large achievable Stokes shifts, narrow symmetric, size tunable emission spectra, and exceptional resistance to photo- and chemical degradation, making them attractive reagents for the long-term visualization of cellular structures and processes. See references 1-4.\nThe methods employed to date for the intracellular delivery of QDs or other nanoparticles (NPs) can be grouped into three generalized categories based on their physicochemical nature. Passive delivery is a nonspecific process that relies on the inherent physicochemical properties of the QD (surface charge and/or functionalization) to mediate uptake.\nFacilitated delivery utilizes a delivery agent (e.g., a cationic peptide or a polymer) that is covalently attached to or electrostatically complexed with the QDs to specifically induce internalization. Both these techniques, while noninvasive, typically utilize the endocytic pathway which results in encapsulation of the QDs within intracellular endolysosomal vesicles and thus requires further strategies to liberate the sequestered QDs to the cytosol if that is ultimately desired. Examples methods of facilitated delivery include using additional chemicals such as sucrose or chloroquine or adding polymers such as polyethyleneimine during delivery to disrupt the endosomes by osmotic shock: such methods are generally quite toxic. Lastly, active delivery methods such as electroporation and microinjection deliver QDs directly to the cytosol through physical manipulation of the cell. However, these are highly invasive techniques that can often compromise the integrity of cellular structures and substantially reduce cellular viability (see reference 10). Thus, each of the previously described methods for delivery of nanoparticles (including quantum dots) is deficient in some way.\nDescribed herein are improvements in intracellular delivery of nanoparticles."}
{"text": "1. Field of the Invention\nThe invention relates to an automatic transmission for automotive vehicles having multiple-ratios wherein ratio changes between at least one pair of ratios is a swap-shift.\n2. Background Art\nSwap-shift transmissions for automotive vehicle powertrains are disclosed in prior art U.S. Pat. Nos. 6,292,731; 5,722,519; 5,553,694; 5,758,302; 6,370,463; and 6,577,939. Each of these patents discloses a control system for a multiple-ratio automatic transmission for automotive vehicle powertrains. The transmission includes first and second gearsets arranged in series so that the torque output element of the first gearset delivers torque to the torque input element of the second gearset. The first gearset is an overdrive gearset or an auxiliary gearset. The second gearset is a multiple-ratio gearset, which is referred to as the main gearset. In the case of the prior art patents identified above, the multiple-ratio gearset is a so-called Simpson gearset, which comprises a first planetary gear unit and a second planetary gear unit with a common sun gear.\nThe torque input element of the first gearset is connected to the turbine of a hydrokinetic torque converter driven by an engine in the powertrain. The torque output element of the second gearset is connected to vehicle traction wheels through a differential-and-axle assembly in known fashion.\nDuring acceleration of the vehicle, as the engine delivers power to the traction wheels, the overall transmission ratio can progress from an initial high torque multiplication ratio to a low torque multiplication ratio with ratio steps between the high ratio and the low ratio.\nThe transmission is characterized by a so-called swap-shift during upshifts from a second overall forward driving gear ratio to a third overall forward driving gear ratio and from the second overall forward driving gear ratio to a fifth overall forward driving gear ratio. Conversely, a swap-shift occurs during downshifts from the third overall gear ratio to the second overall gear ratio and from the fifth overall gear ratio to the second overall gear ratio.\nDuring a swap-upshift from the second overall gear ratio to the third overall gear ratio, the overdrive gearset must be downshifted while the Simpson gearset is upshifted, the shifts in the gearsets being synchronized or coordinated during the overall shift progression. Conversely, a swap-downshift from the third overall gear ratio to the second overall gear ratio requires a simultaneous upshift of the overdrive gearset and a downshift of the Simpson gearset in a synchronized fashion.\nRatio changes are controlled by a pressure operated friction clutch or brake for each gearset. In the case of a 2-3 swap-upshift from the second overall gear ratio to the third overall gear ratio, a reaction brake on the overdrive gearset must be released as a reaction brake for the Simpson gearset is applied. Conversely, on a 3-2 swap-downshift, a reaction brake for the Simpson gearset must be released in synchronism with the application of a friction brake for the overdrive gearset.\nA swap-upshift from the second overall gear ratio to the fifth overall gear ratio requires release of a reaction brake for the overdrive gearset in synchronism with engagement of a friction clutch for the Simpson gearset. The friction clutch for the Simpson gearset drivably connects together two gear elements of the Simpson gearset so that the Simpson gearset gear ratio, following the upshift, is unity.\nIn the transmission of the present disclosure, the overdrive gearset is a simple planetary gear unit with two gear ratios: a low ratio, which is unity, and a high ratio, which is an overdrive. The Simpson gearset is characterized by three forward drive gear ratios (as well as a reverse ratio). The first forward drive gear ratio has high torque multiplication, and the third forward drive gear ratio is unity. The second gear ratio is approximately midway in the torque ratio progression from the highest torque multiplication ratio to the lowest torque multiplication ratio.\nThe overdrive gearset, in combination with the three Simpson gearset gear ratios, is capable of producing an overall ratio range of six gear ratios, since each of the three Simpson gearset gear ratios can be combined with each of the two overdrive gearset gear ratios.\nPrecise synchronization is required to achieve acceptable shift quality during swap-upshifts and swap-downshifts. This synchronization should be maintained throughout the life of the transmission, notwithstanding the presence of wear of the torque transfer friction elements and changes in environmental conditions, such as temperature, lubricant viscosity changes and changes in coefficients of friction for the clutches and brakes.\nEven small errors in synchronization of the gear ratio changes for the overdrive gearset and the Simpson gearset, as the two gearsets are upshifted and downshifted during a swap-shift, will have a significant adverse effect on the overall shift quality. One of these adverse effects is referred to as a ratio “flare” condition. Another is referred to as a ratio “overshoot” condition. These conditions, as well as other error conditions, can cause significant, perceptible torque disturbances at the torque output shaft for the transmission.\nA ratio “flare” occurs when the overdrive gearset begins its gear ratio progression during a swap-shift before the Simpson gearset begins its gear ratio progression. A ratio “overshoot” occurs if the Simpson gearset shift progression ends before slipping of the friction element for the overdrive gearset is completed."}
{"text": "It is well known to use a focused ion beam (FIB) for exposing conductors on an integrated circuit (IC) to aid in debug, or failure analysis, and editing of the integrated circuit. See for instance U.S. Pat. No. 6,225,626, issued May 1, 2001, to Talbot et al.; U.S. Pat. No. 5,140,164, issued Aug. 18, 1992 to Talbot et al.; and U.S. Pat. No. 5,616,921, issued Apr. 1, 1997 to Talbot et al., all incorporated herein by reference in their entireties. U.S. Pat. No. 6,225,626 discloses methods for exposing a selected feature of an IC such as a selective conductor, from the backside of the IC substrate without disturbing adjacent features of the device, such as the active (semiconductor) regions. The method includes determining a region of the IC in which the selected feature is located; obtaining from the backside of the IC substrate a near IR (infra-red) optical microscope image of the region; aligning the optical microscope image with a coordinate system of a milling (FIB) system; and using structures visible in the near IR microscope image as a guide (also described in U.S. application Ser. No. 10/159,527 by Madhumita Sengupta and Mamta Sinha), operating the milling system to expose the selected feature from the backside of the IC without disturbing adjacent features.\nOther aspects of this approach include forming the trench through the backside of the substrate where the trench may be stepped in cross-section, and milling down most of the way through the substrate from the backside surface thereof to just above the substrate circuit elements. Following this an access hole is milled to to expose ILD0 and the metal interconnections (circuitry) typically formed overlying the principal surface of the substrate. The focused ion beam is then used to cut or reform various portions of these metal layers in order to reconnect the transistors formed within the substrate. This operation requires the local isolation of the routed interconnect prior to its construction by means of the deposition of a dielectric layer, as is described in U.S. Pat. No. 5,357,116 by Talbot et al. issued Oct. 18, 1994. This is used typically for failure analysis, debug, and edit (rework) of ICs. It is most typically used in the research and development and manufacturing engineering stages, but is applicable to the reworking or discretionary rewiring of actual production ICs. This is described in Crawford and PE Kudirka, “Electron Microbeam Testing for Large Microcircuit Arrays” IEEE-Proceedings of the Symposium on Electron, Ion and Laser Beam Technology (1971) 131–140. Credence, assignee of this application and of the above-referenced patents, produces commercially available FIB systems, such as the IDS P3Xd instrument and the IDS OptiFIB instrument. A FIB system typically includes a source of ions, such as gallium ions, suitable magnetic and/or electric lenses for focusing the ion beam, and also a chemical source, injected to enhance milling, such as xenon difluoride (XeF2). The OptiFIB includes the focused ion beam column coaxial with a light microscope for optical observation of the milling process when using near infra-red (as described in US Patent Publication number U.S. 20030102436A1), and for navigation (finding specific locations) across the IC die.\nIt has been recognized by the present inventors that prior to performing the final circuit editing operations (cutting and/or filling in of the metal layers), the edit is generally more successful if the access trench itself which is formed on the backside surface of the substrate is etched precisely so that the remaining silicon thickness is uniform across the trench floor. In most applications the trench does not extend all the way through the substrate; instead, it extends to within, for instance, 2 to 5 μm of the principal surface of the substrate, leaving an intervening thin layer of, e.g., silicon. The reason for this is that in advanced ICs, the transistor density on the principal surface leaves no room for very large trenching operations. The actual circuit editing is generally performed through this thin layer by the FIB. For instance, a 30 KeV gallium ion beam, assisted by a chemistry such as xenon difluoride (XeF2), or alternatively with a chemistry less aggressive to silicon such as chemicals containing iodine, bromine, or chlorine, can mill a small via (cavity) through the trench floor by locally removing the last few microns of silicon, and can then add dielectric and conductor, thereby performing the circuit (metal layer) edit.\nSince this silicon layer is the floor of the trench formed from the backside surface-of the substrate, the present inventors have determined that the floor of the trench should be nearly perfectly flat (smooth) and parallel to the principal surface of the substrate; i.e., it is preferable if during the trenching process the remaining silicon has uniform thickness. A non-uniform trench floor may locally enhance the spontaneous etching, which may result in penetration of the ion beam through the silicon into the device active areas and lead to inaccuracy in the circuit edit. Operations such as the automation of edit steps (see for instance U.S. Pat. No. 6,031,229 by Keckley et al., issued Feb. 29, 2000) would depend on this uniformity. Additionally, non-uniformity affects the aspect ratio of the access hole milled through the remaining silicon, and may ultimately influence the success or yield of the attempted edit. Non-uniformity may occur, by way of example, due to the fast silicon removal rate by the assist chemistry which is typically injected into the FIB milling (etching) region. Often the detected non-uniformity is the result of a tilting of the floor of the trench due to the nature of the FIB chemical gas delivery injector.\nIt is believed that others in the field have not adequately recognized the significance of having such a uniform trench floor or at least of having its non-uniformity determined. This has been found by the present inventors to be a significant shortcoming of FIB trench etching processes. In addition to its importance in circuit editing accuracy, thickness uniformity of the trench floor is also critical for enabling accurate navigation by methods such as IR imaging navigation (as described by Talbot et al in U.S. Pat. No. 6,518,571) or voltage/materials contrast navigation, as described in U.S. patent application Ser. No. 10/789,336. In particular, voltage contrast navigation requires etching quite close to the active regions: therefore uniformity of the trench bottom to the die surface is even more important, so as not to damage the active regions. Potentially, as dimensions continue to decrease, milling down to the actual diffusions may be required to achieve sufficient navigation and editing accuracies, which will place even a more stringent requirement on uniformity of remaining substrate thickness."}
{"text": "The subject matter of the present invention relates to a system adapted to be disposed at the surface of a wellbore for determining mud filtrate volume data as a function of depth in a wellbore from a set of parameterized, processed, and wellbore corrected induction logging array signal data, which data is produced by an array induction tool disposed in the wellbore, and for generating an output record medium which reflects and illustrates the mud filtrate volume data in the form of a log, the log being analyzed in order to determine information regarding wellbore formation properties, such as layering, vertical permeability, and formation damage.\nDuring the drilling of a wellbore, mud pumps introduce mud into the well in order to flush rock chips and other unwanted debris out of the wellbore. The mud is introduced into the wellbore under pressure, the mud pressure being slightly greater than the pressure of a formation traversed by the wellbore thereby preventing a phenomenon known as well blowout. The resultant pressure differential between the mud column pressure and the formation pressure forces mud filtrate into the permeable formation, and solid particles of the mud are deposited on the wellbore wall forming a mudcake. The mudcake usually has a very low permeability and, once developed, considerably reduces the rate of further mud filtrate invasion into the wellbore wall. In a region very close to the wellbore wall, most of the original formation water and some of the hydrocarbons may be flushed away by the mud filtrate. This region is known as the \"flushed zone\", or the \"invaded zone\". If the flushing is complete, the flushed zone pore space contains only mud filtrate; in addition, if the flushed zone was originally hydrocarbon bearing, it would contain only residual hydrocarbons. Further out from the wellbore wall, the displacement of the formation fluids by the mud filtrate is less and less complete thereby resulting in a second region, this second region undergoing a transition from mud filtrate saturation to original formation water saturation. The second region is known as the \"transition zone\". The extent or depth of the flushed and transition zones depends on many parameters, among them being the type and characteristics of the drilling mud, the formation porosity, the formation permeability, the pressure differential, and the time since the formation was first drilled. The undisturbed formation beyond the transition zone is known as the \"uninvaded, virgin, or uncontaminated zone\". In FIGS. 1a-1b, a schematic representation of an invasion and resistivity profile in a water-bearing zone is illustrated. FIG. 1a illustrates a cross section of a wellbore showing the locations of the flushed zone, the transition zone, and the uninvaded zone extending radially from the wellbore wall. FIG. 1b illustrates a radial distribution of formation resistivity extending radially from the wellbore wall, into the flushed zone, into the transition zone, and into the uninvaded zone. In FIGS. 2a-2b, a schematic representation of an invasion and resistivity profile in an oil-bearing zone is illustrated. FIG. 2a illustrates the radial distribution of fluids in the vicinity of the wellbore, oil bearing bed. FIG. 2b illustrates the radial distribution of resistivities for an oil bearing zone, similar to the radial distribution of resistivities for a water bearing zone shown in FIG. 1b. Sometimes, in oil and gas bearing formations, where the mobility of the hydrocarbons is greater than that of the water, because of relative permeability differences, the oil or gas moves away faster than the interstitial water. In this case, there may be formed, between the flushed zone and the uninvaded zone, an \"annular zone or annulus\", shown in FIG. 2b, with a high formation water saturation. Annuli probably occur, to some degree, in most hydrocarbon bearing formations; and their influence on measurements depends on the radial location of the annulus and its severity. However, the existence of these zones (the flushed, transition, annular, and uninvaded zones) influence resistivity log measurements and therefore the accuracy of the resistivity log itself that is presented to a client. The resistivity log is utilized by the client to determine if oil exists in the formation traversed by the wellbore. The client is mainly interested in the true and correct value of Rt, the resistivity (reciprocal of conductivity) of the uninvaded zone, since Rt is the best measure of the possibility of oil existing in the formation. However, the existence of the flushed and transition zones in the formation adjacent the wellbore wall adversely affect a measurement of Rt. Therefore, since large amounts of money may be spent based on the resistivity log presented to the client, it is important to understand the true resistivity of the formation in the flushed and transition zones in order to improve the accuracy of the resistivity log in general. Prior art well logging tools function to log the formation and generate signals, which signals are processed by a well logging truck computer situated at the well surface. The well truck computer produces a resistivity log. For a particular depth in the wellbore, the shape of a resistivity radial profile (hereinafter, resistivity profile), produced by the prior art well tool at the particular depth and extending radially outward from a point at the wellbore wall, is shown in FIG. 3. In FIG. 3, the resistivity profile indicates a flushed zone resistivity \"Rxo\" an uninvaded zone (true) resistivity \"Rt\", and a transition zone resistivity represented by an abrupt step function indicated generally by the diameter of invasion symbol \"di\". This step function transition zone resistivity does not accurately reflect the true resistivity distribution of the transition zone in the wellbore; therefore, the value of the resistivity Rt of the uninvaded zone is also adversely affected. The resistivity of the transition zone does not change abruptly as shown in FIG. 3; rather, it changes gradually as shown in FIG. 1b. Therefore, the resistivity profile generated by the prior art well logging tool was at least partially inaccurate. To correct this deficiency, a new transition zone has been defined. In FIG. 4, a resistivity curve is plotted using a set of invasion parameters, that is, the flushed zone 14 resistivity is Rxo, the uninvaded zone 16 resistivity is Rt, and the transition zone 10 resistivity is a gradual decrease (or increase) from Rxo to Rt. The flushed zone 14 extends a radius r1 from the borehole wall radially into the formation; the uninvaded zone 16 begins at a radius r2 from the borehole wall and extends into the formation. Therefore, the transition zone 10 lies within a region defined between radii r2 and r1. In FIG. 4, the transition zone resistivity changes gradually from Rxo to Rt, and not abruptly as shown in FIG. 3. A true and accurate reading of Rt must be obtained to determine if oil exists in the formation traversed by the borehole. In FIG. 5a, a diagram of depth in a borehole versus radius is illustrated, the diagram depicting the radial extension of the flushed zone, the transition zone, and the uninvaded zone over a plurality of depths in a borehole. For example, in FIG. 5a, the flushed zone extends from the borehole wall to radius r1; the transition zone extends from radius r1 to radius r2, and the uninvaded zone extends from radius r2 radially outward into the formation. Notice that there is a distinct difference between radius r2 and radius r1, indicating that the radial width of the transition zone is more than an abrupt step function, as in FIG. 3. In FIG. 5b, a resistivity log or diagram of depth in a borehole versus resistivity is illustrated. In this example, the lowest reading of resistivity is the uninvaded zone resistivity Rt (the conductivity is the highest). The highest reading of resistivity, relative to the other zones, is the flushed zone resistivity Rxo (the conductivity is the lowest). A complete description of the method and apparatus for deriving the invasion parameters useful for defining the transition zone of FIG. 4 is filed Apr. 16, 1991, entitled \"Method and Apparatus for Determining Parameters of a Transition Zone of a Formation traversed by a Wellbore and Generating a More Accurate Output Record Medium\", the disclosure of which is incorporated by reference into this specification.\nMud filtrate invasion analysis from resistivity logs is commonly attempted by qualitative inspection of the separation between measurement displays representing different depths of investigation. However, general comparisons cannot be based exclusively on radial resistivity differences. As described above, an improvement can be obtained by performing an inversion of resistivity differences through an assumed model to yield a parameter dimensioned in length, and FIG. 4 illustrates a model for a radial resistivity profile originating at the wellbore, incorporating a transition zone between undisturbed rock and rock flushed by drilling fluid invasion. The resulting invasion description in length units may adequately describe invasion in water zones, but may not adequately describe invasion in hydrocarbon zones which are the principle zones of interest. A number of variables should be taken into account. For example, hydrocarbon zones are more complicated because depth based variation in the saturation gradient through the flushed zone/undisturbed zone interface may be confused with changes in invasion character. In addition, variations in drilling mud properties between wells will change the radial resistivity profile, as will differences in formation water properties. Furthermore, depth variation in porosity and hydrocarbon saturation must also be taken into account. Therefore, unless these variables are taken into consideration, an interpretation of the resistivity profile could be wrong.\nTo solve this problem, invasion analysis should be performed in the volume domain in accordance with the present invention. This volume domain mud filtrate invasion analysis normalizes the effect of all these variables and is useful for comparing well to well and between zones within a well for both water and hydrocarbon zones of interest."}
{"text": "1. Field of Invention\nThe present invention relates to a method for the containment of fluids from an aircraft commode.\n2. Description of Related Art\nModern commercial passenger aircraft include at least one and often several lavatories. The commodes installed in such lavatories typically comprise a metal bowl having a base that is in fluid communication with a waste holding tank. All portions of the commode, except for the inner surface of the howl, are typically concealed and covered by a removable rigid shroud, which is supported by a plurality anchors affixed to the walls of the lavatory. Ground service maintenance personnel can remove the shroud to obtain access the commode and other aircraft systems that may be installed near the base of the commode and concealed from view by the shroud. The shroud, which supports a toilet seat and a lid, provides aesthetic and easy-to-clean surfaces that are visible to users of the lavatory.\nLiquid and solid waste deposited by a user into the bowl of the commode is transferred to the waste holding tank using a vacuum system. In addition, a quantity of chemical disinfectant fluid, which is typically blue in color and is thus sometimes referred to in the airline industry as xe2x80x9cblue juice,xe2x80x9d is pumped into contact with the inner sides of the bowl during waste transfer to rinse liquid and solid waste from the inner surface of the bowl to the waste holding tank and to sanitize and disinfect the inner surface of the bowl for subsequent users.\nIn many commercial aircraft, there is no fluid-tight seal between the shroud and the rim of the bowl of the commode to avoid the creation of zones within the passenger cabin that are at different air pressures. Thus, a space or gap exists between the shroud and rim through which fluids can flow or leak. And such leaks do occur with some degree of regularity.\nFor example, during ground servicing operations, waste material must be pumped out of the waste holding tank and a fresh supply of chemical disinfectant fluid must be charged into the system. During such ground servicing operations, an oversupply of chemical disinfectant fluid is sometimes pumped into the system. This fluid can flow into and fill up the commode bowl and then flow between the space between the rim of the commode bowl and the shroud. The fluid flowing out of the bowl and into the aircraft can damage avionics and other systems concealed beneath the shroud, can seep into the luggage compartment and other areas of the aircraft, and can flow across the floor of the lavatory and into the passenger cabin. It takes a significant amount of time to clean up an aircraft that has been oversupplied with chemical disinfectant fluid. The aircraft must be removed from service during this time, which results in significant expense to the airline.\nMust less frequently, fluids can flow out of commodes in aircraft lavatories during flight, such as when the commode becomes clogged with waste or paper or when the aircraft experiences turbulence or other in-flight conditions that can result in the escape of fluid. In some aircraft, such as a DC-9, water from the sink adjacent to the toilet drains into the toilet. If the toilet is not flushed and the water is left running by a user, the water will flow out of the bowl and onto the floor. See, e.g., Adrienne P. Samuels, St. Petersburg Times, Sep. 20, 2002, Airliner\"\"s Lavatory Leaks into Luggage Area. A method for the containment of such fluids is needed.\nThe present invention provides a method for the containment of fluids that flow from a bowl of a commode installed in an aircraft lavatory through a space between a rim of the bowl and a shroud positioned above the rim. The method comprises securing a central portion of an apron against an exterior perimeter of the bowl and fastening a plurality of fastening points formed in an outer peripheral portion of the apron to fixed anchors to thereby extend and retain the outer peripheral portion of the apron away from the exterior perimeter of the bowl. Once installed, the apron is able to contain at least a portion of fluids that flow from the bowl through the space between the rim of the bowl and the shroud positioned above the rim.\nIn the presently most preferred embodiment of the invention, the apron comprises a sheet of flexible multilayer material. A first ply of the multilayer material comprises a fluid impermeable barrier ply and a second ply of the multilayer material comprises an absorbent material. The apron is preferably held against the exterior perimeter of the bowl using a band of aluminum. The fixed anchors supporting the outer peripheral portion of the apron are preferably fastened to the walls of the lavatory.\nThe present invention also provides a kit for use in servicing an aircraft lavatory. The kit comprises an apron and optionally a supply of disposable absorbent material, a disposable floor covering for placement on the floor of the lavatory, and a disposable supplemental lavatory service pad.\nThe foregoing and other features of the invention are hereinafter more fully described and particularly pointed out in the claims, the following description setting forth in detail certain illustrative embodiments of the invention, these being indicative, however, of but a few of the various ways in which the principles of the present invention may be employed."}
{"text": "In one aspect, the present invention relates to a method of conditioning raw tar sand to render it suitable for the separation of bitumen from the tar sand by froth flotation. In another aspect, the present invention relates to a method for processing raw tar sand to produce a bitumen enriched tar sand product more suitable for subsequent extraction of bitumen. In yet another aspect, the present invention relates to upgrading the bitumen content of tar sands through separation by froth flotation of a milled tar sand slurry to produce a concentrated tar sand product. In still a further aspect, the invention relates to a method whereby components of the tar sand slurry may be recovered and recycled for reuse. In yet another aspect, the present invention relates to the preparation of a fluid medium suitable for use in the separation of bitumen from tar sands by froth flotation.\nTar sands (also known as oil sands and bituminous sands) are sands which are impregnated with dense, viscous petroleum. Technically, the material should perhaps be called bituminous sand rather than tar sand since the hydrocarbon is a bitumen (i.e., a carbon disulfide-soluble oil). Herein these sands will be referred to as \"tar sands\". Deposits of tar sands are found in many countries, for example, Canada, the United States, and Venezuela. Tar sand deposits found in the Province of Alberta, Canada are being exploited commercially. In the United States, tar sand deposits are found in the Uinta basin of Utah and Colorado. These tar sands deposits contain billions of tons of tar sands which would be very valuable provided an economical method of separating the bitumen from the sand was available.\nAt present, the Canadian tar sands have been the subject of most of the research and development efforts conducted to date. The Canadian tar sands deposits are characterized by a sandstone grain covered by a layer of water which is covered with bitumen. Two methods of recovery of bitumen from these Canadian deposits are known as \"the cold water method\" and \"the hot water method\".\nThe \"cold water method\" does not involve heating of the tar sand other than whatever heating may occur during normal factory operations. The process involves mixing the tar sand with water at ambient temperature, soda ash and an organic solvent. The mixture is then permitted to settle at ambient temperature and a mixture of solvent and bitumen dissolved in the organic solvent rises to the top of the settling zone, much in the manner of cream rising in milk, and is recovered. The cold water method has not achieved commercial acceptance in Canada due to the physical nature of the Canadian tar sands.\nThe best known and most widely used method in the recovery of bitumen from Canadian tar sands is often referred to as the \"hot water method\". In general, the hot water process involves heating the tar sand with steam or hot water, milling the tar sand with a small portion of water at about 175.degree. F., transferring the tar sand into a turbulent stream of circulating water and carrying it through a separation cell maintained at an elevated temperature of about 180.degree. F. In the separation cell, entrained air causes the oil to rise to the top and form a froth rich in bitumen which is then drawn off. The sand settles to the bottom and may be removed. It is believed that separation of bitumen from Canadian tar sands is effected by disruption of the sheath of water surrounding the sandstone grain thereby also disrupting the bitumen covering of the water sheath, allowing the bitumen to rise to the top of the separation cell. The basic hot water system may be enhanced as described in U.S. Pat. No. 3,573,196 issued to Cymbalisty on Mar. 30, 1971, by adding a light hydrocarbon solvent to the slurry which is at a temperature below the boiling point of the solvent and subsequently heating the slurry to the boiling point of the hydrocarbon solvent, which volatilizes the solvent and assists in the separation of the bitumen.\nThese methods have been used to process the raw tar sand as mined. Raw tar sand ores typically contain from about 6% to about 12% bitumen by weight. Thus, treatment of the raw ore by both methods requires enormous quantities of water in relation to the bitumen recovered because of the low bitumen content in the raw tar sands. The cold water method described above as well as using large quantities of water requires large settling ponds. The hot water recovery method permits an energy recovery in the range of 50 to 60% of the bitumen processed. In other words, approximately one-half of the energy value of the recovered bitumen by the hot water method is consumed in its production.\nThe present invention provides a method whereby the raw tar sand ore may be concentrated, thereby increasing the bitumen content of the tar sand processed for bitumen recovery. By utilizing cold (ambient) water in the concentration of the tar sands in accordance with the method of the present invention, the amount of energy necessary to extract the bitumen from the concentrated tar sand is less per unit of bitumen recovery. Another advantage of the present invention is that the froth flotation apparatus utilized is less expensive to purchase, operate and maintain, and consumes less energy than the specially designed attritioning devices utilized in the hot water method. Potentially, an energy recovery rate of 90% is possible. In still another aspect, the invention allows the utilization of standard equipment readily available in the market. The method allows for recycle of the water agents thereby reducing water consumption, which is an important factor in recovery of bitumen from tar sands in the Western United States."}
{"text": "Integrated circuit (IC) chips are often electrically connected by wires (e.g., gold or aluminum wires) to a leadframe or a substrate in a packaging assembly to provide external signal exchange. Such wires are typically wire bonded to bond pads formed on an IC chip using thermal compression and/or ultrasonic vibration. A wire bonding process exerts thermal and mechanical stresses on a bond pad and on the underlying layers and structure below the bond pad. The bond pad structure needs to be able to sustain these stresses to ensure a good bonding of the wire and to prevent damage to the IC chip.\nAlternative bond pad structures were fabricated from the bottom to the top layers, which did not allow metal wiring circuitry and semiconductor devices to pass under or be located below the bond pad structure. For a more efficient use of chip area or to reduce the chip size, it is desirable to form semiconductor devices and metal wiring circuitry under the bond pad. This is sometimes referred to as bond over active circuits (BOAC) or circuits under pad (CUP). At the same time, many processes now use low-k and ultra low-k dielectric materials for the intermetal dielectric (IMD) layers to reduce RC delay and parasitic capacitances. The general trend in IMD designs is that the dielectric constant (k) tends to decrease from the top downward toward the substrate. However, as the dielectric constant (k) decreases, typically the strength of the dielectric material decreases (as a general rule). Hence, many low-k dielectric materials are highly susceptible to cracking or lack strength needed to withstand some mechanical processes (e.g., wire bonding, CMP). Not only do low-k dielectric materials tend to be relatively weak in compression strength, they also tend to be weak in adhesion strength.\nDuring a typical wire bonding procedure, the bond pad structure must withstand compressive and lateral shear forces during the ball squishing stage of wire bonding. These forces may cause cracking in the relatively weak low-k dielectric layers. The bond pad structure must also withstand pulling and torsional forces while the wire is being pulled from the wire bonding tool relative to the bond pad on the chip. These forces may cause peeling or dislocation of the bond pad structure. Hence, there is a need for a bond pad structure that can sustain and better disperse the stresses exerted on it by a wire bonding process, that is compatible with the use of low-k dielectric materials for intermetal dielectric layers, and that will also allow circuitry and devices to be formed under the bond pad."}
{"text": "1. Field of the Invention\nThe present invention relates to a cascade Raman laser.\n2. Description of Related Art\nFIG. 9 is a block diagram showing one exemplary structure of a conventional cascade Raman laser outputting laser light with a wavelength of 1480 nm. (see S. G. Grubb, et al., “High-Power 1.48 μm Cascaded Raman Laser in Germanosilicate Fibers,” in Optical Amplifiers and Their Applications (1995), paper SaA4 for example). This cascade Raman laser 2000 has a structure in which a clad pump fiber laser (CPFL) 100 outputting laser light with a wavelength of 1117 nm as pumping light and a cascade Raman resonator (CRR) 200 that generates laser light with the wavelength of 1480 nm by using the laser light with the wavelength of 1117 nm are connected in cascade. It is noted that marks ‘x’ in the figure indicate fused connecting points. It is also noted that optical fibers represented by thin solid lines have an outer diameter of 125 μm and optical fibers represented by thick solid lines have an outer diameter of 250 μm.\nThe structure of the cascade Raman laser 2000 will be explained concretely from a back of the CPFL 100. Firstly, the CPFL 100 has a structure in which an extension optical fiber 101, a front pumping TFB (Tapered Fiber Bundle) 102, i.e., a multiplexer for multiplexing the pumping light, a highly reflective fiber Bragg grating (FBG, referred to as HR (High Reflection) hereinafter) 103 having reflectivity of about 100% for the wavelength of 1117 nm, a clad pump fiber (CPF) 104 having a double clad structure and Yb is doped in its core, a back pumping TFB 105 and an outputting fiber Bragg grating (FBG, referred to as an OC (Output Coupler) hereinafter) having reflectivity of about 30 to 60% for the wavelength of 1117 nm are connected in cascade by means of fusion splicing. The front pumping TFB 102 is connected with eighteen semiconductor pumping lasers 107 with a wavelength of 989 nm through eighteen pumping optical fibers 108. The back pumping TFB 105 is connected with eighteen semiconductor pumping lasers 111 through eighteen pumping optical fibers 112.\nThe CPFL 100 outputs laser light with the wavelength of 1117 nm by an amplifying action of the CPF 104 to which the pumping light is supplied and by an optical resonator composed of the HR 103 and the OC 106.\nMeanwhile, the CRR 200 has a structure in which an input-side reflector (referred to as a CRRin hereinafter) 201, a Raman optical fiber 202 for Raman amplification, an output-side reflector (referred to as a CRRout hereinafter) 203 and an outputting optical fiber 204 are connected in cascade by means of fusion splicing.\nThe CRRin 201 is composed of five FBGs that selectively reflect lights of mutually different wavelengths. Reflection center wavelengths of the respective FBGs are, from an input side, about 1480 nm, 1390 nm, 1310 nm, 1239 nm and 1175 nm. Meanwhile, the CRRout 203 is composed of six FBGs that selectively reflect lights of mutually different wavelengths. Reflection center wavelengths of the respective FBGs are, from an input side, 1480 nm, 1175 nm, 1239 nm, 1310 nm, 1390 nm and 1117 nm. It is noted that reflectivity of the reflection center wavelengths of the respective FBGs of the CRRin 201 and the CRRout 203 is from 5 to 30% for the FBGs whose reflection center wavelength is 1480 nm and is about 100% for the other FGBs. (It is also noted that in order of the reflection center wavelengths of the respective FBGs of CRRin 201 and CRRout 203 is accordance with U.S. Pat. No. 5,815,518 titled “Article comprising a cascaded raman fiber laser” and is set up to get the high-efficiency optical power with the wavelength of 1480 nm within CRR 200. It is able to operate that in order of the FBGs is not same above-mentioned CRR. No matter what the FBGs is in different orders.)\nNext, operations of the CRR 200 will be explained. When the laser light with the wavelength of 1117 nm is inputted to the Raman optical fiber 202 from the CPFL 100, Raman scattering light with a wavelength of 1175 nm corresponding to a first Stokes wavelength of Raman scattering (referred to as a first Stokes ray hereinafter) is generated and is Raman-amplified. The amplified first Stokes ray is enhanced in terms of its power by multiple reflection of the optical resonator composed of the CRRin 201 and the CRRout 203, functions as a pumping light in due course and generates a second Stokes ray with a wavelength of 1240 nm. Third to fifth Stokes rays of the wavelengths of 1310 nm, 1396 nm and 1480 nm are generated sequentially by the similar operation. Here, the reflectivity of the FBG reflecting the light with the wavelength of 1480 nm corresponding to the fifth Stokes ray is low, so that this light with the wavelength of 1480 nm outputs to an outside of the CRRout 203 through the outputting optical fiber 204. It is noted that because the CRRout 203 has the FBG whose reflecting wavelength is 1117 nm, the laser light with the wavelength of 1117 nm outputted out of the CPFL 100 is blocked from outputting to the outside of the CRR 200 and is efficiently utilized within the Raman optical fiber 202.\nHowever, the cascade Raman laser 2000 having the conventional structure has a problem that there is a case when an output of the cascade Raman laser 2000 decreases even if an output of the semiconductor pumping laser of the CPFL 100 is increased to enhance power of output light."}
{"text": "MicroRNAs (miRNAs) are small, generally between 18 and 24 residues, polyribonucleotides derived from longer hairpin noncoding transcripts in eukaryotes miRNAs play a significant role in cellular developmental and differentiation pathways. Consequently, there have been considerable efforts made to understand and characterize the temporal, spatial and cellular expression levels and patterns of expression of miRNAs to ascertain their precise role in cellular development and differentiation in both normal and disease states.\nmiRNAs are currently studied by, first, obtaining the total RNA from a sample. Next, the total RNA is fractionated into subpopulations by gel electrophoresis or by chromatographic fractionation and size selective elution. Then, the appropriate section of the gel is cut, and the 18-24 RNAs are eluted from the gel, or the eluted fraction containing single stranded RNAs in the size range of 18-24 ribonucleotides is collected. Next, the RNAs are isolated by precipitation and the miRNAs are characterized.\nDisadvantageously, however, these methods do not work well when the amount of sample is small, such as samples from tumor tissue or biopsy material. Further, characterization of the miRNAs isolated by present methods usually comprises a several step amplification procedure followed by detection, quantitation, cloning and sequencing. Because of the large number of steps in these processes and the notorious inefficiencies associated with the repeated purification, the isolation and identification of miRNAs using present methods is time consuming, relatively expensive, requires relatively large amounts of material and is not fully representative of the population of miRNAs expressed within a small sample, such as within a biopsy of a tumor. Additionally, the present methods are not specific to isolating and identifying an miRNA, and therefore, often isolate and identify siRNA, tRNA, 5S/5.8SrRNA and degraded RNA from additional cellular RNAs.\nTherefore, there is the need for an improved method for isolation and identification of miRNAs that is not associated with these disadvantages."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of fault detection systems for automotive vehicle alternator battery electrical energy circuits which operate to continuously monitor the performance of the circuit when in use.\n2. Description of the Prior Art\nVoltage regulator controlled automotive battery charging systems are commonly used to keep an automotive storage battery at full charge level and to provide sufficient electrical power for the operation of the various automobile electrical accessories. Normally an alternator is driven by means of a belt-and-pulley arrangement coupled to the internal combustion engine of the automobile to provide a rectified output current. The alternator operation is controlled by a voltage regulator which senses the battery voltage and controls the field coil excitation for the alternator in response to the sensed voltage. An indicator lamp is typically used in such circuits to provide an indication to the automobile operator of the absence of an alternator output to alert the operator of the automobile to a malfunctioning in the electrical system.\nIn some prior art systems, the DC current is passed through the indicator lamp prior to commencement of the rotation of the alternator rotor to provide the initial alternator field coil excitation. This is the most common method where the steel-pole pieces in the alternator have been designed to retain little or no residual or permanent magnetism and it is, therefore, necessary to pass electrical energy through the field winding to establish the magnetic flux in the presence of the rotor in order to initiate an electrical output. Typically, such prior art systems provide an indicator light warning to the automobile operator only to indicate an alternator malfunction resulting from a low alternator charging output signal. Examples of such prior art systems may be found in U.S. Pat. No. 3,365,646 and U.S. Pat. No. 3,138,751, both of which are assigned to the same assignee as the present invention.\nOther prior art systems have recognized the need for providing a visual indication of other possible faults in the alternator, voltage-regulator and battery system. U.S. Pat. No. 3,673,588 assigned to the same assignee as the present invention is an example of such a system and teaches the use of separate sensing circuits to separately excite different display lamps to indicate either high or low voltage conditions.\nAnother prior art system provides for detection of some fault conditions in an operating alternator battery system. The occurrence of three symptoms primarly due to alternator malfunction are monitored: high voltage, low voltage and excessive ripple. Three separate detectors operate in parallel to turn on the operator indicator light whenever a fixed signal threshold level has been reached for any of the detectors. This system cannot classify or differentiate among the symptoms nor does it provide any protection for the electrical system upon a detection of a fault.\nAdditional prior systems show fault indicator devices for testing the operation of an alternator, voltage-regulator, battery system by the attachment of an external test circuit device. As is common with all such test equipment, it is designed not to be permanently attached to the alternator-battery-voltage regulator system but rather coupled selectively only for the purpose of diagnosing a malfunction. Normally some of the features of the battery charging system are disabled while the automobile engine is started and various electrical loads drawn to simulate normal operation of the alternator-battery system. With various conductor attachments such test equipment uses either separate indicator lights or a common indicator light or designate different types of generalized malfunctions in the alternator-battery system. Various automobile manufacturers' test manuals show such test equipment and the selective disabling of portions of an automobile alternator-battery system. There is usually a description of manually attaching various cables at predetermined positions in the electrical circuit to monitor the performance of an individual feature under a simulated load. In all cases, the alternator-battery system has been disabled in some manner and the operation is merely a simulated test for that portion to which the test device has been connected.\nIn another prior external testing system, the ignition switch is disconnected from the circuit with the engine running and a test device is connected to monitor the alternator output. Disconnection of the ignition switch disables the automotive circuit by removing both the battery and accessory electrical loads. A single fault indicator light is used in a constant on mode to indicate a general type of fault which is characterized by voltage spikes and in a flashing mode to indicate an abnormal voltage output. This system merely detects generalized fault characteristics without distinguishing or identifying the fault.\nAll of the prior art systems use either a single light to indicate a malfunction or a series of display lights coupled to various separate circuits which are individually connected to identify a malfunction. None of these prior art systems has both the advantage of detecting the occurrence of a fault in an operating alternator, voltage-regulator battery system and diagnosing the fault to determine which of the several major electrical component systems may be undergoing a malfunctioning condition. The present invention is intended to be permanently attached to a functioning alternator, voltage-regulator battery system and, at any given time, to indicate not only the fact that a fault has occurred but to identify the major electrical component such as alternator which is undergoing the malfunction."}
{"text": "The present invention relates generally to processes for coating articles by physical vapor deposition and, more particularly to processes utilizing electron bombardment energy.\nSurface coatings have become particularly important in applications where the physical and chemical properties of materials must meet stringent requirements. One such application is the use of protective coatings to increase the corrosion and erosion resistance of interior parts of gas turbine engines. A number of special alloys have been developed for use as coating materials.\nSurface coatings have been applied in a variety of ways. One method useful in coating with high melting temperature materials is physical vapor deposition in which material evaporates from a molten pool in vacuum and travels to the substrate where it deposits.\nIf the coating process conditions are such that vapor particles encounter few other atoms or molecules, then the vapor travels outward from the source along straight lines. The directional distribution of the vapor is described by a well-known Cosine Law which is applicable when the diameter d of the molten pool of the vapor source is small compared to the mean free path .lambda. in the vapor between the source and substrate. The result is that a coating of uniform thickness would be deposited on the interior surface of a sphere tangent to the pool and having a diameter large compared to d.\nIt is well known that deviations from the Cosine Law are observed at high evaporation rates in vapor sources which are heated by an electron beam in a high vacuum, such as 0.1 mT (millitorr) or less. In such cases, the relative deposition rate directly over the source is greater than that given by the Cosine Law. Departures from the Cosine Law have also been observed during high vacuum evaporation of material from resistance heated boats (see K. R. Carson et al, J. Vac. Sci. Tech., Vol. 7, No. 2, pp. 347-350, 1970). Similar departures from the Cosine Law have been observed in molecular beams flowing out of an orifice into high vacuum. Theories indicate that departure from the Cosine Law can be expected when the ratio .lambda./d becomes less than about one.\nHigh rate electron beam evaporation has also been done in chambers where the gas pressures were in the range up to around 30 mT. (see U.S. Pat. No. 3,756,193 to D. C. Carmichael et al.) Collisions between the vapor and the gas result in deviations from the Cosine Law in that evaporated material is deposited on areas which are not in line-of-sight of the vapor source.\nIn many evaporation processes, the substrate subtends only a small solid angle at the vapor source and much of the evaporant is wasted because it is deposited on the walls of the chamber. It is desirable to increase the efficiency of such a coating process by collimating the vapor so that a higher fraction is incident on the substrate. Vapor collimation is particularly important when the evaporant material is scarce or expensive. Such collimation also allows higher coating rates or lower evaporation powers and thus reduces the consumption of energy in the coating process.\nOne vapor collimation method was described by Blecherman et al in U.S. Pat. No. 3,620,815. There the vapor was collimated by surrounding the vapor source with a gas manifold which directed a high velocity stream of inert gas toward the substrate."}
{"text": "1. Field of the Invention\nThe present invention relates to a system for an intelligent attitude and orbit control of a satellite, and more particularly to improved intelligent control apparatus and method where autonomous attitude and orbit control are performed onboard a satellite in normal and contingency modes.\n2. Description of the Conventional Art\nConventionally, attitude determination and control are performed automatically onboard a satellite. There are patents to improved the accuracy of the attitude control system using inertial sensor and the star sensor. See U.S. Pat. No. 4,617,634, entitled, \"Artificial Satellite Attitude Control System,\" dated Oct. 14, 1986. On the other hand, conventional orbit control method requires tracking of the satellite from the ground, calculation of the parameters to change the orbit, and transmitting commands to the satellite. But this method requires much financial and human resources. Also there exists potential for the human error in the process. Thus, Wertz proposed the method of maintaining an assigned orbit without control or intervention from the ground. See U.S. Pat. No. 5,528,502, entitled \"Satellite Orbit Maintenance System,\" dated Jan. 18, 1996. The present invention provides a method and apparatus to control both attitude and orbit of a satellite autonomously based on intelligent decisions onboard the satellite.\nFor controlling orbit and attitude of a satellite automatically, the position of a satellite has to be determined (Navigation technique), and controlled (Guidance technique).\nNavigation and guidance technique is being actively studied and some of them are being registered as patents.\nNavigation technique is reported in the paper, \"Autonomous Navigation\" C. Jayles, F. Alby, J. Berthios, D. Pradines, Spaceflight Dynamics Part II, Edited by Jean-Pierre Carrou, CNES, 1995. And Navigation technique is disclosed as a patent, \"Method and Apparatus for predicting the position of a satellite in a satellite based navigation systems (U.S. application Ser. No. 5,430,657).\nAlso, Guidance technique is reported in \"TOPEX/POSEIDON Autonomous Maneuver Experiment (TAME) Design and Implementation\" by T. Kia, J. Mellstrom, A. Klumpp, T. Munson, and P Vaze, Advances in the Astronautical Sciences, Guidance and Control 1997, Edited by Robert Culp and Stuart Wiens, American Astronautical Socienty Publication, pp. 41-56.\nThere are several reports and registered patents, for attitude control e.g., U.S. Pat. Nos. 5,534,965, 5,458,300 and 5,412,574.\nIt is desirable to have a system where the attitude the the orbit are maintained autonomously onboard the satellite autonomously onboard the satellite during the normal operational mode. Furthermore, in the contingency situation such as when the collision danger exists or when a sensor or an actuator fails, we require the satellite to maneuver automatically to avoid the collision or to operate without the failed sensor. Thus, an intelligent decision making process for the attitude and orbit control are highly desired quality."}
{"text": "1. Technical Field\nThis disclosure relates to electronic design automation (EDA). More specifically, this disclosure relates to creating and using a netlist abstraction.\n2. Related Art\nAdvances in process technology and a practically unlimited appetite for consumer electronics have fueled a rapid increase in the size and complexity of integrated circuit (IC) designs. The performance of EDA tools is very important because it reduces the time to market for IC designs. Floorplanning is an important stage in an EDA design flow that involves determining locations for various objects (blocks, modules, circuit elements, etc.) in the IC layout. The quality of the floorplan can significantly impact the overall quality of the final IC layout.\nUnfortunately, due to the rapid increase in the size and complexity of IC designs, conventional floorplanning tools can take a very long time to generate floorplans and/or generate poor quality floorplans."}
{"text": "The field of the disclosure relates generally to data communications and more specifically, to an optical fiber serial interface module that interfaces between a terminal controller and a data bus.\nAt least some known applications include an ARINC 629 data bus that uses metal twisted pair electrical bus cables, stub cables, bus terminators and current mode couplers (CMC) mounted on heavy metallic panels. However, because of such components, the data bus is bulky, heavy and expensive. Optical communications solutions, such as those that utilize optical fiber as a communications media, are desirable due to the reduced weight, which may be advantageous in an aircraft.\nOne existing solution for implementing an optical fiber data bus incorporates glass optical fibers (GOFs). This system utilizes 850 nm wavelength transmitters and receivers that are packaged individually in a pair, called a Fiber Optic Serial Interface Module (FOSIM). The FOSIM transmitter and receiver have interface electronics to the terminal controller which transmit and receive electrical signal to and from the FOSIM in Manchester bi-phase format. In the typical aircraft application, these FOSIMs are located inside the various avionics subsystems of the aircraft that utilize the data bus for communications. Often, such avionic subsystems are referred to as Line Replaceable Units (LRUs). Inside the LRU, the FOSIMs are mounted along with the terminal controller on a multilayer 6U (full size) VME circuit card.\nHowever, GOFs may be relatively fragile and break relatively easily during installation on a vehicle, such as an airplane. Further, GOFs have a relatively small diameter, which may make optical alignment difficult. Therefore, components associated with GOFs, such as connectors and optoelectronic devices, may be relatively expensive. There is a strong desire in the aircraft production community to develop an optical data bus that uses more robust optical fiber, such as a plastic optical fiber data bus, to replace the current electrical ARINC 629 data bus for future upgrades of such aircraft, although implementations would not be limited to aircraft applications."}
{"text": "Respirators are designed for different types of breathing gas sources depending on their field of use. Stationary respirators in hospitals are usually designed for being connected to a central gas supply system. If such a gas supply is not available, pressurized gas cylinders are used. Respirators that are designed mainly for mobile use make use of a blower, which draws in room air and feeds it directly to the patient while oxygen is possibly added.\nA gas mixer for a respirator intended to be connected to pressurized gas sources is known from DE 24 55 751 B. Two pressurized gas sources are connected in the prior-art gas mixer to a storage tank via two feed lines and corresponding dispensing valves. The gas mixture is sent from the storage tank to the patient by means of a mixed gas line, which originates from the storage tank and likewise has a dispensing valve. To set a predetermined gas mixture, the dispensing valve is at first opened in one of the feed lines until the outlet pressure in the storage tank has increased to a first predetermined pressure value. The dispensing valve is then closed and the dispensing valve in the second feed line is opened until a second pressure value is reached. A defined volume of gas mixture is then removed via the mixed gas line until the outlet pressure in the storage tank has again become established. The prior-art gas mixture is suitable for use as long as a sufficient pressurized gas supply pressure can be made available.\nDE 197 08 094 C2 discloses a respirator, in which ambient air is drawn in with a blower to supply the patient with gas and is compressed until the inspiration pressure necessary for respirating a patient becomes established. Oxygen can be added into a storage tank, which receives the ambient air drawn in, so that different oxygen concentrations can be set in the breathing gas. No provisions are made for operation with pressurized gases alone from a pressurized gas supply system in the case of this respirator."}
{"text": "Power lawnmowers of the rotary blade type include a housing within which at least one cutting blade is mounted for rotation. Commonly, a grass collector is provided, and a chute connects the housing to the collector; typically for a garden tractor arrangement, the chute may have a length of about one meter, and in the event that the collector is not emptied when full, the chute becomes blocked and is difficult to clear.\nVarious proposals have been made for the provision of indicators for the fullness of the collector. Thus, in accordance with one proposal, the collector is supported on one or more springs which deflect under the weight of the load. This proposal utilizes a normally OFF momentary contact switch, which contacts close under a predetermined load, and which serve to cut-off the ignition system of the mower.\nIn accordance with other proposals, optical or sound transducers or electrical resistance sensors are used which change the state of a switch responsive to blockage at the outlet of the chute, and the changed output is indicated either by the operation of an alarm or by changing the operation of the lawnmower, for example.\nOne of the problems encountered with the above proposals is that of transient conditions which trigger spurious signals. This problem was recognized in at least one of the prior art proposals wherein an electronic delay circuit was incorporated to overcome the problem. However, this delay circuit necessitated a delay reset circuit, to overcome the problem of accumulated delays. Other problems are related to the widely differing responses that may be encountered under different conditions such as wetness and length of grass cuttings, and the tendency of the grass cuttings to mat onto all surfaces with which they come into contact and to thereby change the operating characteristics of the transducing elements and eventually to render them inoperative unless carefully and frequently cleaned.\nIt is an object of this invention to provide in a lawnmowing apparatus a simple and robust detection means that is responsive to the state of fullness of the grass collector.\nIt is another object of this invention to provide in such apparatus detection means that has a relatively delayed response time whereby it is inherently less likely to be affected by transient conditions.\nIt is yet another object of this invention to provide in such apparatus, detector means which is disposed in a position in which it is relatively shielded and unaffected by grass cutting build-up and contamination.\nIt is still another object of this invention to provide in such apparatus, detector means is integrated into the operating system of the mower."}
{"text": "This invention relates generally to gas turbine engines and more particularly to rotors or turbine wheels. Still more particularly, the invention relates to a ceramic rotor and a method of forming the same.\nHeretofore, gas turbine rotors or wheels have been constructed of metal. In some instances the metal is forged or cast and then machined into final form. In other instances the wheel or rotor may be fabricated of sheet metal, the material in all cases being required to resist high temperatures and corrosive actions of the gases of combustion. In addition, the rotors must be designed to resist deformation and/or breakage due to high stresses encountered during operation at high speeds.\nPrior rotors have been satisfactory to a certain degree, but in order to secure greater efficiency and other objectives it is proposed to operate the engines at higher temperatures which decrease the active life of the metal rotors causing premature destruction of the engines.\nTo avoid the above objection, it has been proposed to construct rotors of ceramic material which will withstand the high temperatures and corrosive action of the gases of combustion. Such attempts have not been completely successful, however, as the rotors produced do not possess the strength required to resist centrifugal forces due to high-speed operation. Furthermore, the methods employed in producing the ceramic rotors make their cost prohibitive."}
{"text": "Technological Field\nThe present disclosure relates to a technique for estimating the basis weight of a sheet being conveyed in an image forming apparatus.\nDescription of the Related Art\nImage forming apparatuses such as multi-functional peripherals (MFPs) have been widely used. There are many types of sheets available for MFPs. To print appropriately irrespective of the sheet type, control parameters such as sheet feeding speed, transfer voltage, and fixing temperature are set in accordance with the sheet characteristics. Printing without appropriate settings of the sheet characteristics may deteriorate print quality and the like. As such, various techniques for accurately detecting the sheet characteristics have been developed.\nIn relation to this technique, JP 2007-107976 A discloses a sheet type detecting apparatus “capable of detecting the sheet type”. JP 2002-5610 A discloses a thickness detecting apparatus that is “mechanically simple and inexpensive and capable of detecting the thickness of a medium irrespective of the dielectric or optical properties of the medium without the need for maintenance”. JP 5-99605 A discloses a non-contact thickness detecting apparatus that “is not influenced by variation in a gap between a detection electrode connected to an electrostatic sensor and an object to be detected”. JP 2015-193473 A discloses an apparatus for detecting the thickness of the sheet while “removing the influence of the eccentricity of a roller as much as possible”. JP 2010-189137 A discloses an image forming apparatus “capable of determining the sheet type with the sheet stored in a sheet feed tray”. JP 2006-259589 A discloses an image forming apparatus that “has solved implementation issues”.\nThe basis weight of the sheet is an example of the sheet characteristics. The basis weight of the sheet represents the weight of the sheet per area. The unit of the basis weight of the sheet is expressed as “g/m2”. Since the basis weight of the sheet correlates with the thickness of the sheet, the basis weight of the sheet is used in the meaning similar to the thickness of the sheet.\nThe settings of the basis weight of the sheet influence the printing accuracy in particular. Printing without appropriate settings of the basis weight of the sheet may cause various issues such as jams, transfer failures, and fixing failures. Therefore, the basis weight of the sheet needs to be detected accurately. As a method for estimating the basis weight of the sheet, an estimation method using a capacitance sensor has been proposed. This estimation method is the application of the principle that the capacitance of the sheet correlates with the basis weight of the sheet.\nThe capacitance sensor is susceptible to the surrounding environment (e.g., temperature, humidity, and objects nearby). Accordingly, in a case where the basis weight of the sheet is estimated from the capacitance detected by the capacitance sensor, the estimation result of the basis weight varies depending on the surrounding environment. Detecting the basis weight of the sheet while reducing the influence of the surrounding environmental variations is not disclosed in JP 2007-107976 A, JP 2002-5610 A, JP 5-99605 A, JP 2015-193473 A, JP 2010-189137 A, or JP 2006-259589 A. Therefore, there is a need for a technique for accurately estimating the basis weight of the sheet irrespective of the surrounding environment."}
{"text": "The present invention relates to a satellite based communications system for transmitting medically related images from one or more remote locations to a central station and more particularly, to a system and method for converting radiological digital images into a radio signal for transmission to a central station, then at the central station converting the signals back into video or printed images for analysis, all within a limited time.\nThe current clinical practice for mammography depends on a batch-mode diagnostic procedure. This batch-mode operation consists of end-of-the-day collection of data representing each patient's examination at the screening site and then transferring them via surface transportation to the radiologist's location to be loaded on the mammography viewer for the start of the detection and diagnosis process. Typically, it takes less than five minutes for the radiologist to view all images from a single patient exam before deciding whether to view past year examinations and request a magnification view of suspicious areas. The batch approach inherently does not take advantage of the five minute turnaround in which a radiologist can interpret an examination and provide a diagnosis or request that more images be taken. If the images could be evaluated as they are produced, the reader could respond with the examination results while the patient is still at the examination site. As a result, the current problem, in that about 40 percent of patients do not return for follow-up exams, would be solved because the reduced delay between image collection and interpretation so that the need for further imaging is identified before the patient departs the facility. There exists a need to conduct remote radiological exams and to provide results of the exam while the patient is still available for additional exams.\nDigital radiological images typically contain a high volume of data bits. To accomplish remote imaging the high volume of digitized radiological image data must be accurately transmitted via satellite within a short period of time. Systems for transmitting, via satellite, medically related images from remote stations to a central station are not adapted to handle the high volume of data needed for applications such as remote mammography examination. For example, mammography requires high resolution imagery, usually to within a 50 micron diameter. When converted to digital images, such high resolution mammography films require digitization of about 256 million bits per image. Consequently, the ability to transmit these high volume digitized images is limited because most present day communication protocols are designed for low volume data transmissions.\nThe data transmission limitations of the present system also impede the possibility of interactive communications between the remote station and the central diagnostic site. By way of illustration, high resolution mammography films require digitization of about 256 million bits per image and breast mammography typically requires four images per patient; correspondingly, one giga-bit of image data is transmitted per patient. Present systems for transmission of medical data cannot transmit this volume of information within, for example, a 10.5 minute window deemed necessary for interactive communications between the remote site and the professional located at the central diagnostic site.\nThe American College of Radiology and National Electrical Manufacturers Association developed a communications protocol for digital radiological images called Digital Imaging and Communications in Medicine (DICOM). DICOM is a specification of vendor-independent data formats and data transfer services for digital medical images. It was designed to ensure compatibility at the data level between the various manufactures of equipment for handling digitized images. This standard is used extensively in the industry.\nThe Extended Transmission Control Protocol/Internet Protocol, Extended (TCP/IP), was developed by the Network Working Group, of the Internet Engineering Task Force, and published in May 1992 by V. Jacobson, R. Braden, and D. Borman. It is an Internet based communications protocol designed to tolerate unreliable sub networks. This standard is used in the industry.\nIt is thus desirable to provide a satellite communication system utilizing a standard digital format such as DICOM and the Extended TCP/IP protocol so as to enable a remote station to provide high resolution medical images to a central diagnostic station within a time-frame to provide an effective interactive medical exam."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a FINFET transistor structure and a method for forming various FINFET transistor structures. In particular, the present invention is directed to a method for forming various FINFET transistor structures by different possibilities of formation of oxide layers to reduce the leak current of the FINFET transistor structure or to maintain the heat-dissipating of the FINFET transistor structure.\n2. Description of the Prior Art\nOne of the purposes of the development of the semiconductor industry is to enhance the efficiency of the semiconductor devices and to reduce the energy consumption of the semiconductor devices. When it comes to enhancing the efficiency of the semiconductor devices, it is already known in the prior art that different lattice structures may facilitate the mobility of the electrons or the holes.\nFor example, a higher carrier mobility can be observed when a metal-oxide-semiconductor (MOS) is constructed on an n-channel of a (100) lattice of Si, and similarly a higher carrier mobility can be observed when a metal-oxide-semiconductor (MOS) is constructed on a P-channel of a (110) lattice of Si. As a result, when a planar complementary MOS is constructed, Si of different lattices is formed together to get a substrate so that MOS of n-channel is constructed on a (100) lattice, and MOS of P-channel is constructed on a (110) lattice to get a better performance.\nHowever, as the critical dimension of the devices shrinks, in particular for the generations after 65 nm, the multi-gate devices such as a fin field effect transistor (FinFET) is proposed to replace the planar complementary MOS since it is getting harder and harder to reduce the physical dimension of the conventional planar complementary MOS. However, in one aspect, because some of the bottom of the fin field effect transistor is directly connected to the substrate, inevitable leak current is always a serious problem. In another aspect, should the fin field effect transistor be constructed on an SOI substrate to solve the problem of inevitable leak current, another problem arises because of a higher production cost due to much more expensive SOI substrates.\nGiven the above, a novel method for forming a FINFET transistor structure as well as a novel FINFET transistor structure are still needed to bring a resolution to the dilemma."}
{"text": "This invention relates to cabinet enclosures for electrical apparatus such as electrical switching apparatus which may be used as service entrance equipment or similar apparatus positioned ahead of a watthour meter. Such enclosures are sealed shut by an electrical utility to prevent unauthorized access to electrical power. More particularly, the invention relates to such apparatus wherein the cabinet is provided with an internal barrier overlying the electrical apparatus providing limited access to the apparatus through an open side of the cabinet. An external cabinet door closes off the open side of the cabinet. Both the door and the internal barrier are sealed by the utility to restrict the unauthorized taking of electrical power from the apparatus.\nIt is preferred or desirable that the seal for the internal barrier be visible and inspectable from the exterior of the cabinet indicating that the internal barrier is in place, has not been tampered with and is still sealed without breaking the seals for the door and opening the door to visibly inspect the barrier and its seals."}
{"text": "The present disclosure relates to a method for refreshing the content of a memory cell of a memory arrangement and a memory arrangement for carrying out such a method."}
{"text": "PCBs and other substrates are routinely employed to mount and interconnect electronic components, such as integrated circuits (ICs) and discrete components (e.g., capacitors, inductors, diodes and resistors), to form larger assemblies or whole circuits. A bare, or “unpopulated,” PCB or other substrate is assembled through a manual or automated, “pick and place,” process in which components are picked off a tape or other feeder and placed in specified locations and orientations (called “rotations”) on various land patterns provided thereon. Herein, a land pattern is a group of exposed conductive pads configured to receive and provide mechanical support and electrical connections for a specified electronic component. Through-hole and surface mount electronic components are soldered into place on the land pattern to complete a mechanical and electrical connection to the PCB or other substrate on which the land pattern is located."}
{"text": "Conventionally, there are transport devices that, on a production line, transport a workpiece and cause a machine tool at a transport destination to apply predetermined processing to the workpiece. JP 2011-93032A (FIG. 11) discloses a pallet transport device that transports a plurality of pallets having workpieces placed thereon by circulating the pallets on a track-like path composed of two linear portions and two curved line portions that connect between the linear portions.\nThis transport device includes first and second pallet rails that extend in parallel to each other, first and second pallet moving mechanisms that move the pallets guided to the first and second pallet rails along the first and second pallet rails, and pallet transfer mechanisms that transfer the pallets on a semicircular path from an end portion of the first or second pallet rail to an end portion the second or first pallet rail.\nIn the pallet transport device configured in the foregoing manner, pallets that have been moved by the first or second pallet moving mechanism along the first or second pallet rail are transferred by the pallet transfer mechanisms on a semicircular path from the end portion of the first or second pallet rail to the end portion of the second or first pallet rail. The transferred pallets are moved again by the second or first pallet moving mechanism along the second or first pallet rail. In this way, conventional circulation-type pallet transport devices circulate and transport a plurality of pallets on a track-like path."}
{"text": "1. Field of the Invention:\nThis invention relates to a code detection circuit for removing an influence due to an error code transmitted in a band compression facsimile, etc.\n2. Description of the Prior Art:\nIt has heretofore been proposed to detect and remove an error code transmitted in a band compression facsimile, etc. by means of a code detection circuit which makes use of a parity check or cyclic redundancy check code used for a digital circuit of electronic computers, etc. These circuits, however, are very complex in construction and hence are not suitable for widely used data transmission apparatus such as a facsimile, etc."}
{"text": "Typical emergency response vehicles have many different systems for monitoring and responding to various situations and emergencies. For example, the vehicles are equipped with communications equipment that includes both voice and data generating devices such as radios and computers. This and other electronic equipment (e.g., controls for devices such as light bars) crowd the interior space of the vehicle, which is not designed for this concentration of electronics.\nIt is extremely difficult to equip the vehicles with all of the needed communications, monitoring, and response equipment. Standard commercial vehicles are retrofitted with this equipment through a labor-intensive process. Retro fitting the vehicles is often an iterative process, as new equipment replaces old. Advances in equipment allow first responders to perform their jobs more safely and efficiently. However, each time equipment advances, vehicles must again be retrofitted. Furthermore, when the vehicle is no longer used by emergency services, the equipment must be removed from the vehicle through another costly, labor-intensive process.\nIn addition to systems for detecting and responding to emergencies, vehicles must be equipped with various communications systems. For example, in the United States public safety officials including fire departments, police departments and ambulance services primarily use communications systems that work within the VHF and UHF bands. Conventional land mobile radios operate on these and other frequencies. Cellular networks, which operate in the UHF frequency band, are also used for public safety communications systems for both data and voice communications. More recently, the SHF band, such as the 4.9 GHz band reserved by the United States Federal Communications Commission (FCC), have been included in public safety communications systems. Moreover, within these several frequency bands, there are a number of communications standards, such as the IEEE 802.11 protocol, utilized to transmit data. Many other frequency bands and communication protocols are used by emergency service personnel around the country. In order to ensure reliable communications across public safety agencies, vehicles are often now equipped with still further electronics that enable public safety personnel to communicate over several transmissions protocols and/or frequency bands. All of the radios and communications equipment results in a cluttered environment.\nAs technology evolves and finds applications in the area of public safety, emergency response vehicles increasingly carry more equipment to detect and respond to countless situations and emergencies. Typically, individual systems are installed in the vehicle for each of the tasks aimed at emergency responses. For example, a police vehicle monitors traffic using a radar detector. Cameras mounted in an emergency vehicle gather evidence. Many emergency vehicles have light bars mounted to their roofs. Sirens warn citizens of danger. GPS systems inform a control center of the vehicle's location. Vehicles may contain equipment to detect bio-hazards or chemicals in the event of an industrial spill or terrorist attack. Countless other systems are installed in emergency vehicles based on expected situations. This trend can only be expected to continue.\nEmergency vehicles are often equipped with emergency lighting equipment that draw attention to the vehicles and provide visual warning to citizens. Typically this equipment includes flashing or rotating lights, which generating a considerable amount of electromagnetic noise. Because of the noisy environment and to assist in visibility, the emergency lighting equipment is most often housed in a module commonly called a “light bar” mounted to a roof of the emergency vehicle. Installing the emergency light equipment in a light bar lessens the effect the electromagnetic noise has on the operation of sensitive telecommunications equipment inside the vehicle.\nInstalling in emergency vehicles all of this communications, detection and response equipment is costly and labor intensive. All of it is retrofitted into a vehicle manufactured without any accommodation for this special purpose equipment. Some of the equipment, such as radar units and cameras are typically mounted to the front edge of the interior of the roof such that the radar unit and/or the camera extend downwardly to provide views through the front windshield. Power cables are routed from this equipment to the vehicle's power system through the roof lining and down one of the side posts of the car, separating the front and rear car doors, and then to a controller unit, which is located in the trunk, engine compartment or even under a seat in the interior of the vehicle. Many emergency vehicles are equipped with light bars mounted on the roofs of the vehicles. Power and control cables for the light bars are also fished through the side posts and routed to the trunks of the vehicles or to the engine compartments of the vehicles. These cables are fished through the side pillar of the vehicle separating the front and rear doors. Communications antennas are mounted on the roof and on the trunk. Holes are drilled in the car to attach the antennas. Again, cables are routed to a controller in the trunk of the vehicle. Finally, each piece of equipment is wired to controllers in the vehicle's cabin. There are numerous other systems that are regularly installed in emergency vehicles. As technology advances, new devices must be incorporated into emergency vehicles. This requires taking the vehicle out of service for an extended period of time as older devices are removed from the vehicle and newer devices are installed.\nBy their nature, emergencies often require deployment of more emergency equipment than normally in use at any given time. Communities must determine how best to provide for emergency situations that may require quick deployment of additional equipment. Typically, communities rely on resources from neighboring communities. This strategy works as long as the neighboring communities are close by and not affected by the same emergency. For emergencies that affect large areas, however, relying on neighboring communities to loan their resources is not a workable strategy.\nFor example, neighboring communities may face a common emergency such as a hurricane, a terrorist attack or an earthquake. In these types of emergencies, the effected communities will need additional emergency vehicles that are not available from nearby neighboring communities. Moreover, because of the labor intensive and costly installation process, non-emergency vehicles cannot be quickly converted for emergency use. Furthermore, existing emergency vehicles may not have the best combination of equipment for dealing with a particular disaster. The time-consuming installation process prevents vehicles from being quickly adapted to respond to an emergency condition that the vehicle is otherwise not equipped to handle.\nAfter a vehicle is no longer needed by public safety agencies, it is typically sold in the aftermarket. However, all of the communications systems and emergency equipment must be removed from the vehicle before sale. If the vehicle is to be resold at maximum value, the damage to the vehicle done during the process of retrofitting the emergency equipment must be repaired. For example, any holes drilled into the vehicle during installation of the equipment must be patched. The dashboard most likely needs to be repaired because of holes drilled in it to run wiring, mount devices and control units. All of this repairing is expensive and reduces the resale value of the vehicle, which represents a substantial amount of lost revenue to communities.\nAnother problem facing first responders is the lack of a unified communications network for transmitting voice and data. For example, different police departments responding to the same emergency affecting several communities may use different radios. Furthermore, live video taken from one vehicle at the scene of an emergency is not available to other vehicles responding to the emergency. Current attempts to solve communications problems result in even more equipment and radios being installed into vehicles."}
{"text": "1. Field of the Invention\nThe present invention relates to a release agent for metallic molds used for forming plastic molded products.\n2. Description of the Related Art\nPlastics have excellent properties such as easy processability, high productivity, light weight and relative low procuring costs, so they are used for the parts and structural materials for automobiles, autobicycles, scooters, televisions, radios, audio equipment, washing machines, rice cookers, personal computers, portable telephones, game machinery, building materials, office supplies, stationery, toys, sports goods, sports equipment, agricultural tools and marine tools.\nThey are usually prepared by plastic processing methods such as injection molding, blow molding, compression molding, transfer molding, rotating molding, slush molding, inflation tubular film processes, and extrusion molding.\nAs plastics processing are conducted at high temperature above 200° C., additives, monomers, decomposed materials in the plastics are changed to carbonized materials such as tar, pitch and other colored sticky substances, which are apt to adhere to the surfaces of the screw, barrel, die and metallic mold of the plastic processing machine to cause transfer of the carbonized materials to the molded products.\nAbove mentioned transfer makes the surfaces of molded products dirty and makes the dimensions of the molded product inaccurate. As a result, the molded products do not perform required movement or structural functions. Therefore cleaning the screw, barrel, die and metallic mold of plastic processing machines has been required.\nIn case of production changes from specific colored molded articles to noncolored or different colored molded articles, cleaning the screw, barrel and die has been required to avoid cross contamination caused by residual specifically colored resin compounds.\nHowever, there are some problems associated with conventional methods of cleaning the screw, barrel, die and metallic mold.\nThere are problems associated with releasing plastic molded products from metallic molds.\nMetallic molds used for forming plastic molded products have complicated shapes with fine hollow and convex parts.\nTherefore, plastic molded products such as mechanical parts for electronic devices and automobiles have complicated shapes. Consequently it is difficult to release plastic molded products from metallic molds, because both surfaces have complicated contrasting shape, which causes the plastic mold products to ingress into portions of the metallic mold and become anchored thereto.\nTo avoid above mentioned difficulty of releasing both surfaces, there have been efforts to treat the interface between the metallic mold and the plastic molded product.\nThe method of treating the interface involves using a lubricant as a release agent for metallic molds.\nThe lubricant is usually applied by coating it onto the metallic mold surfaces.\nHowever, when a metallic mold that is coated with a conventional release agent is used for a long time, a hard and uneven heterogeneous layer is formed because of the following reasons.\nOne reason is the accumulation of resin and its additives such as antioxdants, metal deactivaters, nucleating agents, anti-ultraviolet agents, antistatic agents, crosslinking agents, valcanization agents and lubricants.\nThe other reason is the accumulation of decomposition materials of resin and its additives.\nThe uneven shape on the surface of metallic mold caused by the heterogeneous layer is apt to be transferred to the surface of the plastic molded product together with the accumulated contaminants from the surface of the metallic mold which will adhere to the surface of the plastic molded product.\nAbove mentioned transferring of shape makes the dimensions of the plastic molded product inaccurate so that it will not perform required movement or structural functions, and adhering of contaminants makes the surface of the plastic molded product dirty.\nThe hard and uneven heterogeneous layer on the surface of the metallic mold surface adversely affects the releasing property.\nTo remove the contaminants and heterogeneous layer from the surface of the metallic mold, it is common to wash the soiled metallic mold after a predetermined time of use.\nAt a plastic molded product facility, washing is carried out after the troublesome task of taking the soiled metallic mold out from plastic processing machine, and disassembling the plastic mold into pieces, which requires time and labor resulting in heavy loss of efficiency and cost.\nEspecially for complicated or large sized metallic molds, removing, dismantling and washing are tremendous and obstacles to improving the productivity and keeping cost down.\nTo decrease the number of times of metallic mold washing, a release agent has been employed.\nAs the release agent, there have been many known types that such as dimethyl polysiloxane of the non-crosslinking type, paraffin wax, higher fatty acid derivatives, metal soaps, talc, mica, polytetrafluoroethylene of the crosslinking type, and the like.\nHowever, conventional release agents have the following drawbacks.\nIn case employing noncrosslinking type of release agent, the release agent stays in specific places of the metallic mold to affect bad influence to the plastic molded product.\nIn the case of employing a crosslinking type of release agent, the release agent will exhibit an excellent releasing property; however, the release agent is apt to adhere to the surface of a plastic molded product and prevent uniform painting or an adhering treatment, resulting to decrease post-processability.\nIn the case of requiring a sufficient releasing effect to achieve a short shot cycle time, the amount of release agent must be increased, result in the accumulation and degradation of the release agent on the metallic mold. This adversely affects the smooth and gross properties and the degradation of mechanical strength, including tensile, elongate and anti-impact strength.\nFurther, it has become a problem that the release agent itself is decomposed by the high temperature heat transfer from the metallic mold.\nFor instance, dimethyl polysiloxane oil which is a representative release agent gradually decomposes over the temperature of 150° C. and rapidly decomposes over the temperature of 200° C. to form a viscous gel-like material which degrades the releasing property.\nFor improving the heat-resistant property of dimethyl polysiloxane, the use of the amino-group or mercapto-group modified dimethyl polysiloxane has been proposed.\nHowever, it has been discovered that the modified dimethyl polysiloxane generates a bad smell such as ammonia gas or mercaptan and discolors plastic molded products.\nAlso, a release agent using dimethyl polysiloxane must be prepared by dispersing it in the water to form micelles using a surface active agent, because dimethyl polysiloxane itself is not compatible with water.\nA release agent using polytetrafluoroethylene has the drawback that it must be baked onto the surface of metallic mold. Therefore many troublesome repeated steps of baking must be conducted notwithstanding its excellent releasing effect and secondary processing properties.\nIt is desirable that the release agent for metallic molds be prepared as an emulsion type from viewpoint of cost, toxicity, ignition, handling and applicability.\nEmulsion type agents are prepared by a method in which a surface active agent, water and non-water-soluble dimethyl polysiloxane oil or wax oil are agitated together to form micelles which are able to be dispersed in water.\nHowever, above mentioned surface active agents react with ingredients that bleed from plastic molded products and form a strong membrane (coated layer) on the surface of the metallic mold.\nThe membrane (coated layer) is hard and has an uneven heterogeneous shape which is transferred to the surface of plastic molded products causing the production of out of standard plastic molded products.\nOn the other hand, when above mentioned release agent of the crosslinked type that requires baking is used, the baked membrane (coated layer) that exhibits the releasing effect is gradually scraped by the resin composition contacting to the metallic mold during each injection shot.\nAs the above mentioned baked type release agent is gradually scraped during each shot of plastic resin composition into metallic mold, the release agent is removed from the surface of metallic mold after from 10 to 20 repeating shots are preformed.\nThen after 10 to 20 shots are preformed, fresh release agent must be applied onto the surface of metallic mold, which is a very troublesome task.\nFurther the above mentioned baked type release agent has the other problem of containing the reactive functional group such as amino, mercapto, isocyanate or vinyl group which is used for conducting crosslinking and the baking reaction on the surface of metallic mold.\nOwing to the reactive functional group, the surface energy of the release agent layer on the metallic mold becomes very large to an increase in the frictional force between the surface of the plastic molded product and the surface of release agent layer, which causes consumption of the release agent layer by scraping during each injection shot, limiting the life of the release agent layer to a term of from 10 to 20 shots of injection mold material.\nThe short life of the release agent requires more frequent application of the release agent by the baking treatment, which leads to accumulation of unnecessary baked crosslinked release agent on the surface of the metallic mold, which adversely affects shape and contaminant transfer to the plastic molded product.\nAs mentioned above, conventional type release agents have drawbacks such as poor release effect, short durability term, difficult secondary processability and difficult application (baking) onto the metallic mold."}
{"text": "A control device with an acceleration sensor and an error control unit for motor vehicles is described in German Patent No. 39 30 302. This control device includes a first acceleration sensor for detecting longitudinal accelerations and a second acceleration sensor, for example, for detecting transverse accelerations. The control device also includes a facility for recording the signals from the acceleration sensors and for deriving a data value on the basis of the two signals. The data value is compared to a preset comparison value in a further facility and an error in one of the sensor signals is detected. The data value is derived, for example, by adding up the acceleration signals or determined as the root of the sum of the squares of the acceleration signals.\nWith this monitoring method, the signals of both acceleration sensors are always used directly for sensor monitoring. This means, first of all, that an acceleration sensor cannot be monitored only on the basis of its own signal. In addition, there is no indication of how to prepare the acceleration quantity prior to monitoring so that it will be largely independent of driving maneuvers performed with the vehicle. In other words, it is not possible to monitor the sensor independently of vehicle movement."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a fin diode structure and a method of manufacturing the same, and more particularly, to a fin diode structure with globally-doped regions in the substrate which is compatible with the process flow of normal fin field effect transistors (FinFET).\n2. Description of the Prior Art\nThe use of fin field effect transistors (FinFETs) in the semiconductor technology keeps rising as the size of technology decreases. FinFETs are advantageous in smaller technologies because of their relatively higher drive current when compared to devices of similar size and because of their general ability to prevent short-channel effects. FinFETs generally have increased drive currents because the gate wraps around the channel such that the effective width of the channel is increased. The increased channel width allows for a greater drive current. Furthermore, by having the gate wrap around the channel, the gate can suppress leakage current through the channel more easily, thereby decreasing the short channel effects.\nThe above-identified advantages of FinFETs have led to their use in smaller technologies, particularly in 32 nm node and smaller. However, the trade-off for smaller size results an increased susceptibility of failure of the FinFET devices due to the electrostatic discharge (ESD) issue. It is well known in the semiconductor field that extremely high voltages can be produced in the vicinity of an integrated circuit due to the build-up of static charges. A high potential may be generated at an input or output buffer of the integrated circuit, which may be caused by a person touching a package pin that is in electrical contact with the input or output buffer. When the electrostatic charges are discharged, a high current is produced at the package nodes of the integrated circuit, and this issue is referred as electrostatic discharge (ESD). ESD is a serious problem for semiconductor devices since it has the potential to damaging the entire integrated circuit. Especially for the FinFET device, the active area width of a FinFET is much smaller than that of another device of corresponding technology size. The smaller width may lead to increased current density in the FinFET when the ESD event occurs, which means that the tolerable and allowable threshold current density is smaller for the FinFET device.\nFor example, FinFETs typically have a threshold current density of 0.1 mA/μm before device breakdown occurs as compared to approximately 2 mA/μm for planar bulk MOSFETs or approximately 1.4 mA/μm for planar SOI MOSFETs. This extremely small current density may cause the dielectric gate oxide to breakdown easily between the active area and the gate and short circuit the gate and the active area. Thus, FinFETs are generally more susceptible to device failures from electrostatic discharge issue because of their relatively small channel width, and a solution is needed to overcome this problem."}
{"text": "Aging of the skin is a complex phenomenon resulting from the interaction of several intrinsic and extrinsic factors. Intrinsic aging is an inevitable, genetically programmed process. Among extrinsic influences (e.g., wind, heat, cigarette smoke, chemicals, etc.), ultraviolet radiation appears to be the single most important factor associated with aging of the skin. The effect of ultraviolet radiation on elastic tissues results in elastosis, which is the accumulation of damaged elastin, resulting in reduced elasticity and resilience.\nElastin is a critical component of extracellular matrix, and is especially abundant in tissues subject to physical deformations, such as lungs, blood vessels and skin.\nThe effect of intrinsic aging on tissue elasticity of mucosal tissues (such as vaginal, oral, or rectal mucosal tissues) and of viscero-elastic tissues (that are lining body cavities such as the respiratory track, the gastrointestinal track, the urinal and bladder track, or the reproductive track) is very similar to the effect of intrinsic skin aging. Elastin fiber production in these tissues is reduced with aging, resulting in reduced responsiveness to stimuli. In the oral cavity, such changes can contribute to a decrease in the health of the gums (leading to reduced resistance to the pressure of food processing), increased gum bleeding, loose teeth, and a general decrease in the visual health parameters of the oral cavity. In the vagina, reduced elastin fiber production could result in stiffness and reduced sexual function, and uterine prolapse is associated with reduced elasticity of the female reproductive system. Reduced elasticity of the bladder can result in urine incontinence. Reduced elasticity of vessel walls can lead to vessel breakage and bruising. In the eye, degenerative changes in elastin fibers in Brunch's membrane can be responsible for deposition of drusen and macular degeneration.\nConsequently, the reduction in elasticity of these tissues results in reduced quality of life and self esteem. Thus, it is desired to have a treatment that can prevent, retard, or reverse the intrinsic and extrinsic aging effects on tissue elasticity.\nTriglycerides are a main constituent of vegetable oil and animal fats, and they play an important role in metabolism as energy sources. However, high triglyceride levels may be associated with a higher risk for atherosclerosis, heart disease, and stroke. (Forrester, J. S., Curr. Opin. Cardiology 2001, 16: 261-264). High triglyceride levels can also increase the risk of thrombosis, which can lead to myocardial infarction (Miller G. J., Atherosclerosis, 2005, 179:213-27). Hypertriglyceridemia is also a well known cause of acute pancreatitis, which can have life-threatening complications (Bae J. H. et al., Korean J Gastroenterol. 2005, 46:475-80). Current approaches for lowering triglycerides include diet and pharmacological agents, such as fibric acid derivatives, fish-oil, and CoA reductase inhibitors (Jonkers, I., et al., Am. J. Cardiovasc. Drugs 2001, 1:455-466).\nUric acid is an end product of purine metabolism. Purines are building blocks of RNA and DNA. Most uric acid produced in the body is excreted by the kidneys. An overproduction of uric acid occurs when there is excessive breakdown of cells, which contain purines, or an inability of the kidneys to excrete uric acid.\nHyperuricemia can play a role in the development of gout as well as many degenerative diseases, such as the Metabolic syndrome, which has been linked to a number of coronary heart diseases and increased mortality (Lee, M-Sh., et al., J. Clin. Nutr. 2005, 14:285:292). Hyperuricemia is also involved in the tumor lysis syndrome (TLS), which is a life-threatening constellation of metabolic derangements arising as a consequence of the release of intracellular metabolites by tumor cells as they undergo necrosis (Zeh, H J et al., J Immunother. 2005; 28:1-9). Uric acid and triglycerides were both found to be positively associated with C-reactive protein (CRP) levels (Garcia-Lorda P., et al., International Journal of Obesity (2005) 1-7).\nThus, it is desired to have a treatment that can prevent, retard, or reverse the negative cardiovascular effects induced by high blood levels of triglycerides and uric acid.\nMalvaceae is a family of flowering plants that includes the mallows, cotton plants, okra plants, hibiscus, baobab trees, and balsa trees. The family traditionally consists of about 1,500 species in 75 genera. Malva sylvestris is a species from the Malva (mallow) genera. The leaves of Malva sylvestris, otherwise known as blue mallow, are rich in mucilage. The mucilage of M. sylvestris is made up of high molecular weight acidic polysaccharides (Classen B, et al., Planta Med 64(7): 640-44 (1988)). The leaf tea is traditionally believed to be useful as an anti-inflammatory, decongestant, humectant, expectorant, and laxative. It has also been used internally for soothing sore throats, laryngitis, tonsillitis, coughs, dryness of the lungs, and digestive upsets. Mallow is also used as a poultice for healing wounds and skin inflammations. In traditional medicine, mallow leaf tea is also used against abnormal growths of the stomach and to alleviate urinary infections (Bisset N G (ed). Malvae folium—Mallow leaf. In Herbal Drugs and Phyto-pharmaceuticals (1994, CRC Press, Stuttgart, pp 313-316). Studies on irritated mucus membranes have shown that the mucilage of Malva sylvestris binds to buccal membranes and other mucus membranes of the body (Schmidgall J, et al. Planta Med 66(1): 48-53(2000)).\nCotinus coggygria extract is traditionally believed to be useful as an anti-microbial treatment, used in the form of external washes. See, e.g., US Patent Applications Nos. 2002/0132021 where the extract is mentioned to be active against E. coli, Staphylococcus aureus and S. cerevisiae, as well as having anti-cancer activity. The dried leaf and twig of Cotinus coggygria is used in Chinese traditional medicine to eliminate “dampness” and “heat”, and as an antipyretic (Huang K. C., The Pharmacology of Chinese Herbs (CRS Press, 1999, pp 193-194). A yellow/orange dye can be obtained from the root and stem and can be used for fabric dying. The leaves and bark are a good source of tannins (Grieve M. A Modern Herbal. Dover Publications, Inc. NY, 1971, pp 779-781).\nThe present invention relates to the unexpected discovery that Malva sylvestris and Cotinus coggygria extracts, when ingested, are both effective for enhancing the elasticity of the skin, urogenital, blood vessel walls, and mucosal tissues, as well as reducing triglyceride and uric acid levels."}
{"text": "This application claims benefit of Ser. No. TO2011A000483, filed 3 Jun. 2011 in Italy and which application is incorporated herein by reference. To the extent appropriate, a claim of priority is made to the above disclosed application.\nThe present invention relates to an electric propulsion system for vehicles, comprising a first and a second electric motor, a gearbox with at least three forward gears and control means arranged to control the first and the second electric motor and to control the engagement of the gears of the gearbox.\nAn electric propulsion system for vehicles of the above-identified type is known from DE-A-101 11 137."}
{"text": "The present invention relates to a novel method and structure for constructing an elevated structure for an artificial landscape. More particularly, the invention is directed to a method and structure which can be used to form a model subroadbed above a base surface so that low-lying areas can more easily be constructed.\nRailroading enthusiasts and hobbyists of all ages have long enjoyed the challenge of model railroading. One of the challenges faced by these hobbyists is constructing a realistic layout that accurately simulates an actual landscape. Before the layout can be constructed, it must first be designed. Designing the layout includes determining the scale, size and overall shape, as well as the time period to be modeled. Further, the modeler must decide what types of industries will be represented on the layout, whether a town will be included, as well as what natural formations, such as trees, lakes and mountains will be present. Certain limitations, such as the available space and the expense involved are, of course, considered when making the above decisions. Further, the layout will include a pattern for the track on which the train will travel. This pattern may involve elevational changes for the track, to simulate grades, bridges and tunnels. The layout may also include low-lying areas, to simulate such things as rivers, ditches and valleys. After the layout is designed, it must then be constructed.\nIn general, railroad transportation involves a locomotive that pulls the rolling stock, which may include passenger cars and freight cars. The locomotive and the rolling stock are supported and travel along a track that is in turn supported by a roadbed. The roadbed is supported upon a subroadbed structure. Thus, in constructing a model railroad layout it is necessary to construct the subroadbed upon which the track is placed. The subroadbed that is constructed must conform to the grades in the layout, and support the track and roadbed that are placed thereon.\nIn the past, when low-lying areas were to be constructed, a benchwork support system is used. The benchwork is constructed of a series of wooden supports, which support pieces of a base material. Various levels of the layout may be created by supporting the various base pieces with the wooden supports at the needed heights. This allows a low-lying area to be created by supporting different base pieces at different heights. Low-lying areas would include streams or rivers, valleys, ditches and ravines. Basically, it may be desirable to simulate any low-lying area which exists in the real world. However, as may be appreciated, constructing such a benchwork is not a simple task, and requires the use of power and hand tools, as well as a high degree of skill. Further, once the benchwork is constructed, the modeler is somewhat restricted in changing the layout if any changes in the benchwork are required.\nPrevious methods for creating a graded subroadbed for a model landscape have also been difficult, time consuming, and noisy. The needed inclines or declines were typically constructed from wood and required the use of power tools, hammers and nails. The nature of the materials used made it difficult to construct an incline or decline with a uniform and continuous grade. The difficulty increased significantly when an incline or decline was desired to be curved so that a rise or fall in elevation could continue throughout a radius in the layout. Further, the previous methods and devices for constructing a terrain grade and subroadbed resulted in a relatively heavy layout. If the layout was desired to be somewhat portable, the added weight made it more difficult to relocate the layout.\nTherefore, a method and a structure are needed that can be used to quickly and easily create a relatively lightweight subroadbed on an artificial landscape that more easily allows low-lying areas to be created. Still further, a method and structure are needed that allow a modeler to more easily change the overall layout without having to replace the layout base or benchwork."}
{"text": "The thermal therapy (hyperthermia) method is attracting attention as a method for cancer therapy. As for this thermal therapy method, attention is focused upon the fact that a cancer cell or a cancerous tissue is more vulnerable to heat than a healthy cell. By heating the cancer-affected portion for example to 43 degree C. for a certain length of time, the cancer lesion can selectedly be necrotized. In this thermal therapy method, an aqueous sol of magnetic fluid comprising a complex of dextran or its derivative and magnetic iron oxide, dextran magnetite for example, is injected into the affected portion, to which a strong magnetic field is applied from outside to selectively heat the cancer lesion.\nIn such a method for therapy, in order to inductively heat the magnetic material having been injected into the affected portion, it is necessary to irradiate a strong magnetic flux from outside to the affected portion by using a magnetic field generation device. A magnetic field generation device preferable for such an inductive heating comprises a transformer, an inverter circuit connected to the primary side of the transformer and a series resonance circuit including a magnetic field generation coil and a resonance condenser connected to the secondary side of the transformer. The magnetic field generating device transmits through electromagnetic induction an alternate output of the inverter circuit connected to the primary side of the transformer to the series resonance circuit connected to the secondary side of the transformer. An alternate magnetic field is generated by applying a current of high frequency to the magnetic field generation coil. On this occasion, if the frequency of the alternate output transmitted to the series resonance circuit on the secondary side conforms with the resonance frequency of the series resonance circuit, the amount of current passing in the series resonance circuit is large enough to generate an alternate magnetic field of needed intensity at the magnetic field generation coil, making it possible to irradiate from the magnetic field generation coil a strong magnetic flux to the affected portion.\nOn the other hand, if the frequency of the alternate output transmitted to the series resonance circuit on the secondary side does not conform with the resonance frequency of the series resonance circuit, only a small amount of current passes in the series resonance circuit, failing to generate an alternate magnetic field of needed intensity at the magnetic field generation coil. Here, that the frequency of the alternate output transmitted to the series resonance circuit on the secondary side does not conform with the resonance frequency of the series resonance circuit means that the switching frequency of the inverter circuit does not conform with the resonance frequency of the series resonance circuit. In this case, the impedance on the primary side of the transformer becomes large, drastically reducing the electrical power transmitted from the primary side to the secondary side of the transformer as well as the amount of the current passing in the series resonance circuit, so that the intensity of the alternate magnetic field generated from the magnetic field generation coil is drastically decreased.\nInductance of the magnetic field generation coil may vary significantly according to factors like the size and the position of the affected portion. Accordingly, the true resonance frequency cannot be obtained by calculating the resonance frequency only based upon the design values of the magnetic field generation coil and the resonance condenser, which leads to a problem that a magnetic field of the needed intensity cannot be provided to the affected portion in cancer therapy. More importantly, according to various experiments and analyses concerning the series resonance circuit having been carried out by the inventor of the present invention, the inductance of the magnetic field generation coil and the capacity of the resonance condenser may vary from the design values according to environmental conditions such as temperature. A magnetic field of the needed intensity may not be generated during the whole therapy at the magnetic field generation coil if the switching frequency of the drive signal supplied to the inverter circuit is fixed based upon these design values. By thermal therapy, therefore, it is crucially important to make the switching frequency of the drive signal supplied to the inverter circuit follow accurately the resonance frequency of the series resonance circuit when the resonance frequency varies according to the size and the position of the affected portion and changes in the environmental conditions.\nFor the purpose of solving this problem, in the Japanese Patent Laid-open Publication No. 2002-360712 is disclosed a thermal therapy device in which the switching frequency of the drive signal supplied to the inverter circuit is conformed to the resonance frequency of the series resonance circuit. This thermal therapy device comprises a transformer, an inverter circuit connected to the primary side of the transformer, a series resonance circuit connected to the secondary side of the transformer which includes a magnetic field generation coil and a resonance condenser, and a PLL (phase locked loop) controlling unit for supplying a drive signal to the inverter circuit. The PLL controlling unit detects a current passing between the inverter circuit and the transformer (a current on the primary side) and generates a drive signal by PLL controlling in which the phase of the current on the primary side is compared with the phase of the drive signal so that the switching frequency of the drive signal to be provided to the inverter circuit is conformed with the resonance frequency of the series resonance circuit. With such a constitution, when the resonance frequency of the series resonance circuit varies according to changes in environmental conditions or the like, and the current on the primary side decreases, the PLL controlling unit generates the above mentioned drive signal in such a way that the current on the primary side has the same intensity as in the case that the frequency of the current on the secondary side conforms with the resonance frequency, in other words, in such a way that the switching frequency of the drive signal supplied to the inverter circuit conforms with the resonance frequency of the series resonance circuit. A current having the same frequency as the resonance frequency passes in the series resonance circuit. Therefore, if the resonance frequency of the series resonance circuit varies according to changes in environmental conditions or the like, the switching frequency of the drive signal provided to the inverter circuit can follow the resonance frequency."}
{"text": "The present disclosure generally relates to an eye tracker, and specifically relates to an eye tracker that can be embedded in a frame of a compact augmented reality (AR) display.\nHead mounted devices (HMD) for virtual reality (VR) systems typically have a mirror arranged in between a lens and a display. The mirror directs light reflected from the eye to a camera that is used for eye tracking. Because virtual reality systems have relatively large form factors, and have room between the lens and the display, the mirror can be positioned behind the lens.\nAugmented reality systems typically rely on wearable devices that have much smaller form factors than classical VR head mounted devices. For example, augmented reality devices may be near-eye-displays that resemble eyeglasses. In such devices, a mirror cannot be conveniently placed behind a lens. Instead, augmented reality devices often include a camera that images the eye directly. However, to avoid disrupting the view of the user, the camera is typically placed far from the optical axis, e.g., in the frame of the device. This results in a very off-axis view, which hinders accurate tracking of the eye. Furthermore, the camera is exposed, and can easily become dirty, scratched, or damaged."}
{"text": "A tool for preparing coaxial cables with a ringlike outer conductor, can only ensure a fast, reproducible and high quality assembly of connectors, if it meets a number of requirements:\nCutting the central conductor, the dielectric, the outer conductor and the sheath of the cable, in accordance with the requirements of the connector to be fitted thereon.\nCutting cable conductors with no burrs or filings. A thin burr, or filings, inside the outer conductor may generate intermodulations (generally known as PIM, the acronym of Passive Intermodulations).\nAutomatically ensuring preparation accuracy in relation to a reference (a crest or a groove of the outer conductor), regardless of the accuracy wherewith the cable is sawed.\nThe invention particularly pertains to a tool meeting all the aforesaid requirements. The invention also provides a method for preparing coaxial cables with a ringlike outer conductor, as well as a tool specifically designed for that purpose.\nThe invention thus provides a method for preparing coaxial cables which have a sheath, a ringlike outer conductor, a central conductor and a dielectric separating the outer conductor from the central conductor, by cutting the outer conductor, the dielectric and the central conductor to length, characterized in that the cable is cut perpendicularly to its axis, in a position approximately determined in relation to a groove of the outer conductor; in that, with the help of a tool with a blade, the sheath is sheared off by rotation in a helical motion, while a guide of the tool comes to rest on the central conductor and, under the action of a positioning and abutment device of the tool, one or more supplementary blades of the tool, cut the outer conductor, the dielectric and the inner conductor, the inner conductor and the dielectric on the one hand, and the outer conductor on the other hand, respectively, and possibly the outer conductor is treated by means of a flaring tool.\nThe invention thus also provides a tool for preparing coaxial cables which have a sheath, a ringlike outer conductor, a central conductor and a dielectric separating the outer conductor from the central conductor, characterized in that it comprises\na blade for shearing the sheath off;\na guide which may rest on the central conductor;\na positioning and abutment tool;\none or more supplementary blades of the tool, for cutting the outer conductor, the dielectric and the inner conductor, the inner conductor and the dielectric on the one hand, and the outer conductor on the other hand, respectively.\nAccording to the invention the tool may preferably be designed to be provided with a flaring tool for treating the outer conductor.\nAccording to further preferred features of the invention,\nthe tool may in particular be composed of a cable sheath peeling knife, which holds the blade fastened by a screw, said peeling knife being screwed on a guide of a blade holder and said blade holder being fastened by a screw and holding a supplementary blade and an axle about which the inner guide may rotate, secured by a screw;\nthe split rings may in particular be housed between two guides and may rotate freely relative to each other and to a spring and open apart while being held in a coaxial position with respect to their housing by a split zigzag or polygonal flat spring;\nsaid that flat spring has a pentagonal shape;\nblade may in particular be fastened to form an angle of about 6 degree with respect to a plane perpendicular to the axis of the tool, which allows it to move in a helical path, as the tool rotates and advances;\nsaid blade holder may in particular comprise two blades, in order to allow the cutting of the central conductor of the cable, in a way shifted from the cutting plane of the dielectric and of the outer conductor;\na blade holder with three blades may in particular be used;\na flaring tool may in particular be used, composed of a body, an inner guide and a blade fastened by a screw in an aperture, said blade haying a dimension (v) which is greater than the greatest radius of the dielectric and smaller than the greatest inside radius of the outer conductor;\na blade may in particular be used which is made of one piece with the blade holder and in that the cutting edge of the blade is shifted through a distance (*) from the radius.\nAccording to further preferred features of the invention,\nthe method may in particular involve that, with the help of a handle or of a drilling machine, secured at the end of the blade holder, the tool is turned, while it is slightly pushed towards the cable.\nthe method may in particular involve that, while the tool rotates, a guide a rolling movement on the inner surface of the central conductor of the cable, avoiding the formation of filings and reducing the effort to drive the tool."}
{"text": "Auto theft is a common occurrence; alarms and ignition kill switches have been developed, but have not proved adequate against organized car theft. Brazen hold-ups of drivers is an increasingly common form of car stealing. In such circumstances, key-dependent devices are of little use. Such devices are also expensive to very expensive depending on sophistication.\nThere exists accordingly a need for car theft deterrent devices which are effective in the absence of the owner, and despite the thief having the car ignition key. Car tracking systems are being offered but are too costly except for the most expensive of automobiles. But, the great bulk of cars stolen are of only moderate cost, making tracking devices too expensive for use."}
{"text": "In the crude oil market, there is a premium for refineries to process feeds derived from more refractory crudes into high value finished products. These refractory feeds are characterized by a higher density resulting from an increased fraction of aromatic compounds with a low hydrogen content. Therefore, these feeds require deeper hydrogenation (more extensive saturation of aromatic compounds) to produce fuels that meet specifications, such as viscosity, cold flow properties, cetane index, smoke point, emission requirements. The volumetric productivity of such a process provides higher benefits to refineries that deploy hydrogenation (HYD) catalysts powerful enough to meet the required hydrogenation activity.\nAlthough hydrogenation is an important aspect of refining high density feeds, deep hydrogenation catalysts need to deliver more than maximizing the increase in volume between feed and refined products. In the wake of European environmental legislation, there is a global trend towards more stringent legislation that mandates environmentally friendly transportation fuels. An example is the mandate that all diesel produced in and imported into the US will have to be ultra-low-sulfur diesel (ULSD) as of Dec. 1, 2014. This change to transportation fuels with the low sulfur content allows for application of newer emission control technologies that should substantially lower emissions of particulate matter from diesel engines. Environmental mandates to remove contaminants like sulfur from refined oil products implies a need for catalysts which combine a deep hydrogenation (HYD) capability with a deep hydrodesulfurization (HDS) activity. This deep hydrodesulfurization activity is particularly important in refining high-density feed stocks, because a higher density typically implies a higher concentration of contaminants such as organic sulfur and nitrogen molecules, as well as metals and asphaltenes (i.e. multi-ring based aromatic compounds with a low solubility even in the most aromatic solvents).\nDeep hydrodesulfurization is accomplished most efficiently through a combination of: 1) hydrogenation (HYD), which releases sulfur atoms after saturating the ring structure of parent aromatic compounds and 2) hydrogenolysis (HYL), which breaks the bond between a sulfur atom and the carbon atom(s) in the sulfur containing molecule, such as aromatic ring compounds. This implies that optimum catalysts for the deep hydrogenation of high-density feed stocks are expected to exhibit an appropriate balance between hydrogenation and hydrogenolysis functions. Catalysts that are suitable for hydroprocessing (e.g. hydrodesulfurization and hydrodenitrogenation) generally comprise molybdenum or tungsten sulfide or sulfocarbide, in combination with an element such as cobalt, nickel, iron, or a combination of thereof. Attempts have been made to modify the morphology of hydroprocessing catalysts to provide ways to control their activity and selectivity. For example, U.S. Pat. No. 4,528,089 discloses that catalysts prepared from carbon-containing catalyst precursors are more active than catalysts prepared from sulfide precursors without organic groups.\nConsiderable experimental and modeling efforts have been underway to better understand complex metal sulfide catalysts, including factors that control metal sulfide morphology. U.S. Pat. Nos. 7,591,942 and 7,544,632 demonstrate that sulfiding a self-supported multi-metallic catalyst in the presence of a surfactant amine gives a catalyst comprising stacked layers of molybdenum or tungsten sulfide with a reduced number of layers in stacks. A lower number of layers in stacks implies the presence of smaller crystals of molybdenum, tungsten or molybdenum tungsten sulfides, which can result in larger surface areas available for catalysis.\nThe saturation of aromatic compounds in distillate fractions, e.g. vacuum gas oil, heavy coker gas oil or diesel fuel has also drawn attention of researchers. A high aromatic content is associated with high density, poor fuel quality, low cetane numbers of diesel fuel and low smoke point values of jet fuel. High hydrodenitrogenation (HDN) activity is typically associated with aromatic saturation activity of nickel tungsten sulfides catalysts; while high hydrodesulfurization (HDS) activity associates with high hydrogenolysis activity of cobalt molybdenum sulfide catalysts.\nThere is still a need for improved catalysts with improved catalytic activity and resistance towards deactivation, specifically self-supported mixed metal sulfide catalysts for use in the hydroprocessing of lower grade, more refractory hydrocarbon feeds, capable of generating low aromatic products meeting new emission requirements. There is also a need for a better understanding of structure-performance relationships for these catalysts to design next generation of highly active self-supported catalysts for use in hydroprocessing."}
{"text": "1. Field of the Invention\nThe present invention relates to a mattress, particularly an adjustable mattress with supporting force regulated by second flexible members.\n2. Description of the Prior Art\nA case for a mattress's supporting force regulated with second flexible members shown in the patent application Ser. No. 12/877,275, MATTRESS STRUCTURE (hereinafter referred to as cited case 1), submitted on Sep. 8, 2010. Cited case 1 has major features as follows: A mattress structure is a mattress body and comprises: a fixed plate, which has a plurality of flexible accommodating holes and at least one first flexible member; a moving plate, which has at least one second flexible member corresponding to the flexible accommodating hole, the moving plate is below the fixed plate and has a distance with the fixed plate for moving; at least one lifting device, which is combined with the bottom of the moving plate, the moving plate and the second flexible member move up and down via the lifting device in order to let that the second flexible member penetrates through the flexible accommodating holes and protrudes out of the top surface of the fixed plate.\nThe patent application herein which is different from cited case 1 is to combine and accommodate first flexible members and second flexible members in a carrier and each of the first flexible members and each of the second flexible members are encased in packages to become a flexible support unit. In a regular state, one of the first flexible members sustain weight of the bed surface by means of two bearing points in a stretchable direction contacting both the bed surface and the bottom surface; one of the second flexible member allows its bottom to protrude from or be concealed in the carrier's bottom surface regularly. In addition, the bottom surface of the carrier can be a wire netting structure and each of the first flexible members is fastened to the bottom surface by means of fastener components without exposure of the first flexible member due to the shifted package."}
{"text": "Olefin polymerization catalysts are of great use in industry. Hence there is great interest in finding new catalyst compounds, particularly compounds that can be easily supported and or modified by altering the ligands. There also exists the need to reduce catalyst costs by developing catalyst systems that are efficient utilizing simple alkyl aluminums rather than alumoxanes that are traditionally more expensive.\nSilica supported chromium tris(amide)/alumoxane olefin polymerization catalysts are found in U.S. Pat. No. 6,040,399 where calcining at greater than 500° C. is disclosed to reduce polyethylene molecular weight (column 4, line 24). Further the U.S. Pat. No. 6,040,399 examples use Cr(III) tris(trimethylsilyamide) which may or may not react with the support and as noted they found only formation of oligomers if using silica calcined above 500° C. Further in U.S. Pat. No. 6,040,399, comparative example 4, describes a catalyst system with a supported Cr(III) tri(trimethylsilyamide) catalyst activated with triisobutylaluminum yielding an extremely low activity polymerization catalyst.\nU.S. Pat. No. 6,011,127 discloses a process for making polyethylene comprising polymerizing ethylene in the presence of a supported catalyst comprising a chromium compound such as a chromium amide, an alumoxane, and an organometallic alkoxide.\nU.S. Pat. No. 5,723,399 discloses an ethylenic polymerization catalyst comprising a chromium compound, a carrier, alumoxane, and a transition metal compound having a conjugated π (Pi) electron as a ligand, where in the chromium compound is not calcined.\nU.S. Pat. No. 6,528,596, discloses a catalyst system prepared by contacting silica successively with biscyclopentadienylchromium (II), pentafluorophenol and ethylenebis(indenyl)zirconium dichloride without the use of an amide intermediate. The resultant polymers had a narrow polydispersity.\nU.S. Pat. No. 5,723,399 describes polymerization of an olefin using various chromium compounds and a metallocene, but without pentafluorophenol. It provides comparative examples using a metallocene alone and various chromium compounds alone.\nAdditional references that are of interest include: JP-11092522 A2 (1999-Apr.-06); JP-11228620 (1999 Aug. 24); JP-10338707 (1998-Dec.-22); M. Enders et al., Organometallics, 2001, Vol. 20, pages 5005-5007; K. H. Theopold, Eur. J. Inorg. Chem., 1998, pages 15-24; U.S. Pat. Nos. 3,709,853; 2,825,721; 6,011,127; and 6,040,399."}
{"text": "Web materials are known in the art. For example, a polyester nonwoven that is impregnated and/or coated with a detergent composition is known in the art as shown in prior art FIGS. 1 and 2. As shown in prior art FIGS. 1 and 2, a known nonwoven substrate 10 is made of dissolvable fibers 12 wherein the nonwoven substrate 10 is coated and/or impregnated with an additive 14, such as an active agent, rather than the additive 14, such as an active agent, being present in the dissolvable fibers 12. An example of such a web material is commercially available as Purex® Complete 3-in-1 Laundry Sheets from The Dial Corporation.\nFurther, an article of manufacture formed from a cast solution of a detergent composition is also known in the art and is commercially available as Dizolve® Laundry Sheets commercially available from Dizolve Group Corporation.\nHowever, such known web materials and/or articles of manufacture exhibit negatives that make them problematic for consumers. For example, the known web materials and/or articles of manufacture are relatively stiff and/or inflexible as measured by the Plate Stiffness Test Method described herein. Further, the web materials and/or articles of manufacture typically deliver such a low level of detergent composition and/or detergent actives that the cleaning performance is less than desired by consumers. Another negative with is that the web materials and/or articles of manufacture may leave remnants of the web material and/or articles of manufacture after the washing operation, for example the polyester nonwoven substrate does not dissolve during the washing operation.\nIn light of the foregoing, it is clear that there is a need for a web material that overcomes the negatives associated with known web materials and/or articles of manufacture described above."}
{"text": "1. Field of the Invention\nThe present invention pertains to an exercise device for a human user; more particularly, to an exercise device that activates, strengthens, and conditions core and other muscles while the user performs sit-ups and other core exercises.\n2. Discussion of the Related Art\nExercise equipment comprises many different devices for exercising the core muscles, including the back and abdominal muscles of the human body. The exercises may be used as physical therapy to strengthen muscles in an area that has been injured or suffers from some other pathological condition. In fact, exercises that invoke the transversus abdominus muscles (corset muscles) and other core muscles are used to treat and manage 1) spondylosis (degenerative osteoarthritis of the joints between the spinal vertebrae) and spondylolithesis (displacement of the vertebrae or vertebral column, sometimes called curvature of the back); 2) chronic or acute lower back dysfunction (LBD); 3) spina bifida; 4) hamstring strains; and 5) problems with the adductor magnus. In addition to physical therapy for rehabilitation, core exercises may be used to strengthen the muscles thereby improving the user's athletic performance or the user's appearance by reducing sagging and otherwise tightening and improving the tone of the core muscles.\nPrior art core exercise devices, such as a slant board or roman chair, usually adopt “sit-up,” “crunch,” or “leg-lift,” like body movements or combinations of two or more of these body movements. The goal of such exercise equipment is to strengthen core muscles including the abdomen and back muscles. Specific names of such muscles include, but are not limited to, the rectus abdominis, the external and internal abdominal obliques, adductor magnus, transversus abdominis, multifidus, and erector spinae.\nIn addition, exercise devices have been created which require a user to maintain balance and equilibrium. The goal of exercising with these types of equipment is to improve a user's core muscles or “stability muscles,” thereby also improving the user's balance. Such muscles are located in a person's lower-torso region. Examples of equipment that facilitate balance are inflatable exercise balls such as the Swiss ball and the BOSU® ball, cylindrical rollers such as Air Roller™, and balancing pads such as DuraDisk™.\nHowever, what is needed is a core exercise device that facilitates exercises dedicated to strengthening the aforementioned core muscles while improving balance. The present invention combines exercises for strengthening the core muscles with a balancing exercise without having to perform separate exercises on separate exercise devices."}
{"text": "In an effort to raise the level of calendering, the only practical possibility has been to increase the number of calendering nips. This has led to a more complicated calender construction and more difficult control and threading of the paper web. Particularly, in reference to on-line machines, a solution must be found to contradictions arising from high running speed and full-speed web threading. Attempts have been made to solve these problems with various belt and shoe calenders for extending a calendering nip to enhance thereby the nip operation. For example, paper to be calendered with belt calenders is carried by means of an endless belt to a preliminary contact with a hot calender roll, thus enabling the creating of a steep temperature gradient favourable from the viewpoint of calendering. By means of the belt, the effective length of a nip is extended by virtue of the preliminary contact as well as for the reason that the belt material enables the use of polymers which are considerably softer than those used in roll coatings, without problems resulting from deformations related to heat. With a nip more extended than in a supercalender or soft calender, the press impulse applied to the paper can be increased without the pressure peak becoming excessively high and without reducing the bulk.\nOne solution for a belt calender has been disclosed earlier e.g. in Finnish patent publication No. 95061. A calender embodiment according to this cited publication is shown schematically in FIG. 1 of the drawings, depicting the prior art.\nIn the prior art calender shown in FIG. 1, a calendering nip N is formed between a heatable hard roll 1 and a calendering belt, preferably a metal belt 5, supported by an elastic surface roll 2. The metal belt 5 is an endless belt and could be e.g. steel in its material. The belt has been extended over the nip roll 2 and a reversing roll 4, the former being lined with an elastic coating 3. In this prior art calender, the calendering nip N is formed between the heatable hard roll 1 and said metal belt 5 supported by the calender roll 2 lined with an elastic coating. Such a solution is largely consistent with the nip of a soft calender, wherein, however, the metal belt 5 enables both sides of a paper W to be subjected to a substantially identical treatment and, thus, the glazing to occur simultaneously on either side of the paper W."}
{"text": "The semiconductor industry is rapidly decreasing the dimensions and increasing the density of electronic circuitry and electronic components in microelectronic devices, silicon chips, liquid crystal displays, MEMS (Micro Electro Mechanical Systems), printed wiring boards, and the like. The integrated circuits within them are being layered or stacked with constantly decreasing thicknesses of the insulating layer between each circuitry layer and smaller and smaller feature sizes. As the feature sizes have shrunk, patterns have become smaller, and device performance parameters tighter and more robust. As a result, various issues which heretofore could be tolerated, can no longer be tolerated or have become more of an issue due to the smaller feature size.\nIn the production of advanced integrated circuits, to minimize problems associated with the higher density and to optimize performance, both high k and low k insulators, and assorted barrier layer materials have been employed.\nTitanium nitride (TiN) is utilized for semiconductor devices, liquid crystal displays, MEMS (Micro Electro Mechanical Systems), printed wiring boards and the like, and as ground layers and cap layers for precious metal, aluminum (Al) and copper (Cu) wiring. In semiconductor devices, it may be used as a barrier metal, a hard mask, or a gate metal.\nIn the construction of devices for these applications, TiN frequently needs to be etched. In the various types of uses and device environments of TiN, other layers are in contact with or otherwise exposed at the same time as the TiN is etched. Highly selective etching of the TiN in the presence of these other materials (e.g. metal conductors, dielectric, and hard marks) is mandatory for device yield and long life. The etching process for the TiN may be a plasma etching process. However, using a plasma etching process on the TiN layer may cause damage to either or both the gate insulating layer and the semiconductor substrate. In addition, the etching process may remove a portion of the semiconductor substrate by etching the gate insulating layer exposed by the gate electrode. The electrical characteristics of the transistor may be negatively impacted. To avoid such etching damage, additional protective device manufacturing steps may be employed, but at significant cost.\nWet etching methods for TiN are known. Such methods may include use of etchants containing hydrofluoric acid in combination with other reagents. However, the selectivity with silicon based dielectrics and metals (e.g., Al) is not sufficient and other exposed metals in the device may also undergo corrosion or etching.\nHydrogen peroxide/ammonia/EDTA (ethylenediaminetetraacetic acid) mixtures and hydrogen peroxide/phosphate mixtures have been disclosed as means of overcoming the acidic HF based etchants. However, the etch rates obtained are insufficient.\nThus, there is a need for TiN etching solutions that have high etch rates, but have low etch and corrosion rates for other semiconductor materials which are exposed or in contact with the TiN during the etching process."}
{"text": "Conventional economy brand caulking guns are generally manual trigger-operated devices incorporating a unidirectional gripping assembly which urges a piston rod forward to eject the compound from a cartridge. As shown in FIG. 1 (prior art), a common gripping assembly includes a piston rod 4 passing through a gripping plate 6 that, in turn, is controlled by the operation of a trigger 8 to grip and move the piston rod 4 in a forward direction for dispensing the caulking composition. The trigger 8 cants the gripping plate 6 clockwise when the trigger is actuated until it grips the piston rod 4, thereafter urging the piston rod 4 forward. Upon release of the trigger 8, the gripping plate 6 is returned to a neutral upright position to allow retraction of the piston rod 4. However, the piston rod 4 also passes through a release plate 10 that prevents undesired retraction of the piston rod 4. The release plate 10 is commonly located at the rear of the gun and remains disengaged from the piston rod 4 during dispensation. However, when the trigger 8 is released, the release plate 10 rides the piston rod 4 rearwardly and becomes increasingly canted relative to the piston rod 4 until it engages the piston rod to thereby prevent inadvertent retraction. The user must depress the release plate 10 to free the piston rod 4 for full retraction.\nOver time the alternating clockwise and counterclockwise torques imposed on the steel thrust components of the above-described mechanism tend to create wear and tear tolerances. Specifically, the hole through the gripping plate 6 and the release plate 10 become enlarged. These wear tolerances combine with the inherent manufacturing tolerances to the point that they effect the operation of the gripping plate and release plate. The effects includes excessive trigger slack and \"piston rod lock\". With regard to FIG. 1, piston rod lock occurs when the hole in the gripping plate becomes so large that the gripping plate tends to return past its normally-neutral upright position (.o slashed.1) to an oppositely-canted position (.o slashed.2) Because of the undesired cant of angle (.o slashed.2), the rod 4 again grips the piston rod 4, and the rod 4 cannot be pulled back even if the release plate 10 is released.\nThere have been previous attempts at implementing wear compensating devices to prevent the above-described problem. For instance, U.S. Pat. No. 5,370,282 to Sedlmeier shows an advancing mechanism (6) with two pivoting clamping levers (12, 13) mounted on the piston rod (7). Clamping lever (12) is for forward advancement while lever (13) prevents retreat. Two springs (14, 15) are required, each in contact with a corresponding clamping lever (12, 13). In addition, an unlocking lever (17) operates via a transmission element 18 to substantially simultaneously unlock both of the clamping levers (12, 13) to free the piston rod (7). The clamping lever 13 rests in a pendular manner in a recess of the frame 5. When the piston rod 7 moves forward, the clamping lever 13 pivots forward along with rod 7 to a perpendicular position (clockwise about the top). When the rod 7 is released, the clamping lever 13 pivots backward (counterclockwise about the top) until it inclines enough to bite the rod 7, thereby locking it in position. When the cartridge is spent, the unlocking lever (17) operates through the transmission element 18 to kick clamping lever 13 forward to release the rod. Thus, the use of the transmission element 18 to unlock both of the clamping levers (12, 13) somewhat resolves the problem of a worn gripping plate that returns past its normally-neutral upright position to improperly regrip the piston rod. However, the '282 transmission element 18 is a small tube that rides the piston rod. It is only capable of pushing against the immediate periphery of the hole in the clamping lever 13. This gives the transmission element 18 very poor leverage against the rod, and as wear becomes worse the transmission element 18 loses its ability to return the clamping lever 13 to a vertical position. The problem of piston rod lock reoccurs. Moreover, the thrust mechanism is complex inasmuch as it requires two separate pivoting clamping levers (one for moving the piston rod forward and one for locking the piston rod in the forward position), two springs each contacting one of the clamping levers, plus a release lever, and the transmission element. The additional parts and labor greatly increases manufacturing costs and assembly time.\nSimilarly, U.S. Pat. No. 5,823,403 to Schneider shows a drive mechanism that provides a double-lever action to bias the gripping plate into a neutral position in a controlled manner to permit the piston rod to be retracted. A gripping (or recoil) plate 52 is normally biased by spring 50 into a canted position to prevent rearward movement of the piston rod 26. However, when the piston rod 26 moves forwardly the recoil plate 52 pivots forward along with rod 26 to a perpendicular position (counter-clockwise about the bottom) and rod 26 becomes free to slide through the recoil plate 52. The recoil plate 52 is situated rearwardly of the frame and this necessitates a backplate 40. The '403 patent appears to be functionally similar to the above-described '282 patent except that the release lever 17 of the '282 gun kicks both levers 12 and 13 out via sleeve 18, while the '403 lever 62 only kicks out the recoil plate 52. Otherwise, both mechanisms requires three separate pivoting levers (one for moving the piston rod forward, one for locking the piston rod in the forward position, and a release plate). Moreover, both require two springs each contacting one of the clamping levers, plus a transmission element. The sleeve 18 only pushes against the immediate periphery of the hole (giving it very poor leverage against the rod), and this provides a temporary solution at best to the wear problem because sleeve 18 eventually loses its ability to prevent piston rod lock.\nIt would be advantageous to provide a more reliable solution to piston rod lock and other wear-and-tear problems to ensure proper advancing and releasing operation in an economical two-lever drive assembly."}
{"text": "Modern electronic devices, such as personal computers, notebook computers, netbook computers, cell phones, smartphones, mobile internet devices, and so on, utilize radio modules to communicate over wireless networks. Often, such devices may include two radio modules, for example, a wireless local area network (WLAN) radio module operating in compliance with an Institute of Electrical and Electronics Engineers (IEEE) 802.11 standard, and a personal area network (PAN) radio module operating in compliance with a Bluetooth protocol. Wireless local area networks are certified by the Wi-Fi Alliance, and are typically referred to as Wi-Fi networks. Alternatively, a wireless wide area network (WWAN), such as a Worldwide Interoperability for Microwave Access (WiMAX) compliant network or a Third Generation (3G) compliant cellular network may be used in combination with a Bluetooth (BT) network. Users of these devices typically will operate both radios simultaneously, for example, by streaming music to the device via the WLAN network and listening to the streamed music with Bluetooth connected headphones.\nDue to channel adjacency, simultaneous operation of WLAN and Bluetooth radio modules may result in collisions on the transmission medium, thereby causing interference and/or packet loss. One of the common approaches to mitigate such issues is by using a media access control (MAC) coordination interface, often referred to as “coexistence interface,” that attempts to coordinate the operation of the two radios in a manner that minimizes collisions by partitioning the medium usage between radios. Current coexistence interfaces are typically designed as point-to-point interfaces to coordinate between two distinct radio technologies. Such designs may differ by the semantics of the signals in the interface, and/or by the behavior of the MAC layer or higher network layer. For example, a Bluetooth radio should behave differently when working with different types of radios, for example using different parameters for voice calls depending on the technology with which it cooperates.\nIncreasingly, mobile devices are incorporating three or more radios into a single device which should be designed to coexist in a similar manner as a two-radio device. As multiple radios are, however, combined into single modules, the pin count for a multiple radio interface may be greater than desired for a practical module. Typically, a two-wire interface may be utilized to implement a coexistence interface between two unique radio modules. Adding a third radio module would then involve a four-wire interface, thereby leading to a higher pin count, especially as additional radio modules are added. Furthermore, routing complexity may increase with additional wires used in the coexistence interface. When trying to reduce the number of pins, it may become a challenge to have different semantics on the wires between the two or more interfaces. As a result, the wires cannot be easily combined using simple logic gates or the like. An additional challenge is controlling the direction of the signals. While a Wi-Fi-BT coexistence interface may have two wires with two signals going in opposite directions, a WiMAX-BT coexistence interface may have two wires going from the WiMAX module to the BT module with two signals going in the same direction. This means if a Wi-Fi and WiMAX combination module were provided, to share a wire that functions in different directions, the correct timing between the Wi-Fi and WiMAX combination module and the BT module would have to be accommodated to ensure that the wire is never being driven at both ends simultaneously which might cause electrical malfunctions and/or render a module non-functional.\nIt will be appreciated that for simplicity and/or clarity of illustration, elements illustrated in the figures have not necessarily been drawn to scale. For example, the dimensions of some of the elements may be exaggerated relative to other elements for clarity. Further, if considered appropriate, reference numerals have been repeated among the figures to indicate corresponding and/or analogous elements."}
{"text": "It is known to use liquid cryogenic fuel in engines and to pre-pressurise the cryogen before injection into the engine. The present invention is a development of the cryogenic engine system described in U.S. Pat. No. 6,983,598 B2. A cryogenic engine is an engine which is driven by the expansion of a cryogenic fuel.\nIn existing systems, such as that described in U.S. Pat. No. 6,983,598 B2, problems have arisen in delivering low quality (near fully liquid phase) cryogen to the engine due to boil off caused by heat transfer from the surroundings into the feed lines. For subcritical feed pressures, this manifests itself as a significant increase in specific volume of the fluid.\nOther issues which have been identified in existing systems include:                a reduction in tank pressure as liquid cryogen is removed; the reduction in tank pressure reduces the feed pressure to the inlet of the cryogenic pumping or injection system, potentially reducing flow rates. This issue needs to be addressed to allow continuous running of the engine, otherwise tank pressure may reduce to a point where cavitation at the pump inlet prevents sufficient liquid flow being delivered to the engine for the desired power output. It is also possible that cavitation could damage pumping equipment;        reducing the pressure under which a cryogenic liquid is stored may also cause unwanted boiling of the liquid. Increased flow rates reduce the unwanted boiling of liquid nitrogen between the tank and the pump, and between the pump and the engine;        flow pulsations in the feed line caused by the cyclic operation of the cryogenic pump and injection valves. Pulsations caused by valve inlets may interact with those caused by the pump, may be amplified by cavitation, and ultimately may affect flow rates into the engine, causing unsteady engine operation;        difficulty in introducing cryogen and heat exchange fluid (HEF) during the short time frame prescribed by the engine operating speed, which can impact the efficiency of operation of the engine; and        the limited period of time in which heat exchange can take place between the cryogen and the HEF, which, again, can impact the efficiency of operation of the engine.        \nUS 2008/0271455 A1 relates to a heat exchange subsystem of a continuous flow cryogenic engine system intended to circumvent the problem experienced in conventional cryogenic engines of icing up from moisture found in ambient air, which can decrease the efficiency of the heat exchangers in their function. Heat is generated by burning hydrocarbon based fuel (e.g. propane). By mixing the heated air and combustion products with liquid cryogenic fuel inside a chamber, the subsystem delivers a continuous flow of high pressure gas at near ambient temperature which can be used in a connected expansion device to extract work. However, the gaseous oxygen required for combustion is provided by the vaporisation of the cryogen within the combustion chamber, but no consideration is given to the flash point of the combustible fuel which dictates whether the fuel will ignite. Moreover, the reduction in temperature and consequent contraction of the combustible gases due to vaporisation and heating of the cryogen mitigate the expansion of the cryogen and, therefore, the build-up of pressure.\nTherefore, there is a need for an improved cryogenic engine system which overcomes these issues."}
{"text": "This invention relates to a metal-oxide-hydrogen cell and more particularly, to a variable conductant heat pipe which is an integral part of such a cell to improve thermal control.\nMetal-oxide-hydrogen batteries are well-known as typified by U.S. Pat. No. 4,115,630. Typical is the nickel-hydrogen cell such as shown in FIGS. 5 and 6 of U.S. Pat. No. 4,115,630. Cells of that type have been used successfully as prime energy storage systems on earth satellites. The cell shown in U.S. Pat. No. 4,115,630 has been used, for example, aboard the NTS-2 (Navigation Technology Satellite-2) spacecraft. In such spacecraft designs battery modules can be located at various places on the spacecraft. For example, a preferred location is on an outside panel that may view cold space. Heat which is generated by the battery will be rejected to space as a result of a completely passive direct radiation system. In such passive systems, radiator sizing can be optimized to maintain battery temperatures at acceptable levels during various cyclic operations, for example, during period of trickle charge. However, passive systems are not generally effective during prolonged time periods where the radiator is exposed to the sun or heat generation rates vary vis-a-vis normal spacecraft orbital operation. For example, during periods of eclipse heat generation can cause significant rises in battery temperature.\nFor example, in the case of the NTS-2 a passive system was designed to maintain the battery at temperature levels of approximately 0.degree. to 5.degree. C. However, during eclipse battery temperatures were actually in the order of 15.degree. C. to 25.degree. C., that is, an increase of approximately 10.degree. C.\nAccordingly, within this technology a need exists to define a system having higher heat rejection rates to flatten the battery temperature profile and thereby lower the average operating temperature of the battery.\nAccordingly, a requirement of thermal control in such spacecraft batteries exists which cannot be adequately satisfied by using completely passive direct radiation systems. Such completely passive systems do not have the ability to provide active thermal control during periods of high battery heating such as during periods of high battery heat dissipations such as eclipse periods.\nIn spacecraft power systems there is a standing design goal to increase usable energy density. Increases in energy density are accomplished by reducing overall weight of such spacecraft cells and/or increasing the depth of discharge of such systems. These design techniques in turn aggrevate problems of thermal control by increasing both the end of discharge and the average operating temperature. Passive systems are simply not satisfactory for such storage systems having very deep discharge cycles. If such were used, unacceptable temperature swings could result during operation of the system. For example, in the case of a low orbit spacecraft application, with a 90 minute orbital period stored energy would have to be removed within 30 minutes yet completely returned within 60 minutes. The limited amount of time for energy dissipation results in higher temperatures. Passive systems do not have the capability for variable heat control.\nIn addition to the passive system as used in a spacecraft environment, the prior art defines other techniques of passive thermal control in context of a non-spacecraft cell designs. One such technique is shown in Sanderson, U.S. Pat. No. 3,498,844, which shows the use of a heating vent 25 in FIG. 2 of that patent, located at the center of a primary fuel cell evaporator 20. In this passive heat transfer system, heat produced at the electrodes 8 and 10 causes the temperature of the working fluid to increase until it is high enough to open a pressure regulating valve 27. As a result, water vapor will be exhausted through vent 25 as a result of the pressure increase in the system.\nAnother passive system is shown in U.S. Pat. No. 3,865,630 which in FIG. 2 also shows a heat pipe 27 utilized as a heating device for electrolyte 23 in a molten salt battery. As shown in that system, the heat pipe is a unitary body which is integrated into the framework of the battery. The heat pipe cannot actively vary its conductants in response to a change in temperature of the system.\nMetal-oxide-hydrogen batteries, in particular, nickel-hydrogen batteries, are assembled utilizing a multitude of cells having an electrode stack contained within a pressure vessel. As shown in FIG. 5 of U.S. Pat. No. 4,115,630 the pressure vessel 90 contains an electrode stack which is axially aligned on a centrally disposed rod 52. Two sets of electrodes comprising a first back-to-back stack of positive electrodes 74 are coupled to a positive bus bar 92 and a second stack of negative electrodes 78 are disposed relative to a negative bus bar 94. At one end of the cell, coupled to the negative bus bar, a negative support 114 is used as a negative terminal. At the other end, a fill tube 126 is disposed coaxially within a positive terminal 112 coupled to the positive bus bar 92. As indicated, such a cell relies on passive thermal control. The predominant transfer path for heat which is generated in the electrode stack within the cell is radially through the hydrogen gas and then through the pressure vessel to the battery structure and finally, radiated to space."}
{"text": "1. Field of the Invention\nThe present invention generally relates to the production techniques of integrated circuits and, more particularly, to a dynamically switched voltage screening method for quality assurance testing of integrated circuit dies on a wafer.\n2. Description of Related Art\nIntegrated circuit (IC) chips or dies fabricated on the same wafer have a wide range of performance in critical parameter characteristics due, in a large part, to process variations. As a result, out-going quality and reliability is compromised by the inconsistent performance of IC chips from the same wafer lot. Typically, stress testing is performed on the wafer during this fabrication process to eliminate the weaker chips from entering the next phase of the production cycle. For example, one of the tests regarded as one of the most severe, is that of subjecting the devices to particularly high temperatures; typically between 85xc2x0 C. and 150xc2x0 C., for accelerating the infant mortality of chips in a wafer. This test, commonly referred to as xe2x80x9cburn-inxe2x80x9d, has the objective of stimulating the failure of those devices which have developed some defects during the fabrication process and/or during handling. However, this test does not segregate parts based solely on performance characteristics. Rather, the stress is uniformly applied to all chips of a semiconductor wafer. Additionally, this test is time consuming, equipment intensive, and costly to perform.\nAlternatively, the industry has administered a voltage stress test for implementing a screen with less degrading effects. A voltage screen is basically a higher than normal voltage applied during wafer test: that effectively causes defects to manifest as failures in IC chips during subsequent verification or operational testing. The problem associated with applying a higher than normal voltage level across a semiconductor die is that some dies processed with short channel lengths have a higher tendency to fail when exposed to higher voltages. These short channel length devices would not otherwise be failures except for their vulnerability to high voltage exposure. The failures are not related to defects, rather, to the over stressing of the short channel lengths in the die. Thus, as applied, the voltage screen can be responsible for false failures, an undesirable quality assurance test result. Nevertheless, the industry standard has been to continuously apply a voltage stress at one voltage level to all die on a semiconductor wafer.\nAs discussed by Lee and Sonoda, in xe2x80x9cTEST SYSTEM FOR NARROWING THE RANGE OF PERFORMANCE CHARACTERISTICS OF MONOLITHIC INTEGRATED CIRCUITSxe2x80x9d, IBM Technical Disclosure Bulletin, Vol. 15, No. 4, September 1972, some circuits, such as those employed with FET technology, have performance characteristics that can differ by as much as 100% due to tolerances in the threshold voltage. Because of these wide differences, there are at one end of the distribution curve chips exhibiting fast response time and high power dissipation, and at the other end of the distribution curve chips having slower circuit response time and lower power dissipation. Thus, sorting the wafer at the IC chip level during quality assurance or reliability testing would be advantageous to identifying the more vulnerable short channel length devices and exposing them to less stress, thus, eliminating false failures during voltage screening.\nIn the prior art, parts have been screened during wafer testing to classify individual IC chips at various speeds. In U.S. Pat. No. 5,196,787 issued to Ovens et al. on Mar. 23, 1993, entitled, xe2x80x9cTEST CIRCUIT FOR SCREENING PARTSxe2x80x9d, a test circuit was developed on the die to measure the DC characteristics of a device, which in turn, enabled one to estimate the AC characteristics. The AC characteristic estimations were then used to screen parts into various speed classes. However, no suggestion is made to dynamically switch or adjust the stress test levels based on the operational parameters measured.\nAnother method for determining the operational speed of an IC chip is disclosed in U.S. Pat. No. 5,099,196 issued to Longwell et al. on Mar. 24, 1992, entitled, xe2x80x9cON-CHIP INTEGRATED CIRCUIT SPEED SELECTION.xe2x80x9d By forming an oscillator in an IC semiconductor chip to generate pulses representative of the speed of other components formed in the chip, the operational speed of the oscillator (typically, a ring oscillator), and therefore, that of the other components formed in the semiconductor chip, can be determined. Again, the stress test levels are not altered in response to the operation speed measurements taken.\nIC chip segregation tests also include bit-pattern recognition on each device under test. This method is particularly useful in memory device testing. In U.S. Pat. No. 4,335,457 issued to Early on June 15, 1982, entitled, xe2x80x9cMETHOD FOR SEMICONDUCTOR MEMORY TESTINGxe2x80x9d, semiconductor memory devices are tested using a special purpose computer which employs simple test patterns to determine the weakest bits of the device, and then tests only these relatively few xe2x80x9cweak bitsxe2x80x9d along with structurally and operationally adjacent bits using highly complex test patterns to determine if the device is functioning properly. Bit pattern recognition, however, is not a stress test screen. Thus, no adjustment of stress test levels, predicated on the bit pattern results, is either taught or suggested by this prior art.\nStill, other methods may be employed to distinguish the IC chips based on variations in the operational parameters. However, independent of the method chosen, some functionally operating IC chips continue to remain vulnerable to excessive stress test screening levels due to chip-to-chip process variations.\nBearing in mind the problems and deficiencies of the prior art, it is therefore an object of the present invention to provide an apparatus and method that determines the speed of IC chips on a semiconductor wafer and adjusts the stress test levels based on the speed measured for each device.\nIt is another object of the present invention to provide an apparatus and method for effectively protecting the short channel IC chip population with a lower voltage during voltage stress testing of a semiconductor wafer.\nA further object of the invention is to increase the outgoing quality and reliability of a semiconductor die using a voltage screen without falsely rejecting the short channel die during voltage stress testing of a semiconductor wafer.\nIt is yet another object of the present invention to provide an apparatus and method for segregating IC chips capable of a higher voltage withstand level without compromising the resultant yield from the wafer lot.\nAnother object of the present invention is to increase the measure of reliability of the devices on a semiconductor wafer by assigning supply current limits as a function of device speed.\nStill other objects of the invention will in part be obvious and will in part be apparent from the specification.\nThe above and other objects and advantages, which will be apparent to one of skill in the art, are achieved in the present invention which is directed to, in a first aspect, a method for testing integrated circuit semiconductor devices comprising the steps of: providing a wafer containing a plurality of integrated semiconductor devices; measuring a desired parameter of the devices; and, applying a stress test to the devices wherein test conditions of the stress test are adjusted based on the desired parameter measurements of the devices. Measuring a desired parameter first comprises verifying functionality of at least some of the integrated semiconductor devices at a set of operating conditions. The method further comprises the steps of: verifying device functionality at nominal operating conditions after the stress test; and, classifying the devices as failed if the devices do not function properly after the stress test.\nThe measuring of a desired parameter comprises measuring the operational speed of the devices prior to applying the stress test.\nAdditionally, applying a stress test comprises applying a first voltage at a value higher than the device normal operating voltage to the devices with a first measured operational speed, and a second voltage at a value lower than the first voltage to the devices with a second measured operational speed, the first operational speed being less than the second operational speed.\nThe method also comprises assigning supply current specification limits and measuring the limits of the devices such that a first supply current specification limit is assigned to devices with a faster of the operational speeds and a second supply current specification limit is assigned to devices with a slower of the operational speeds, the first supply current limit being greater than the second supply current limit.\nThe current invention is directed to, in a second aspect, a method for testing integrated circuit semiconductor devices comprising the steps of: providing a wafer containing a plurality of integrated semiconductor devices; determining functionality of the integrated semiconductor devices at nominal operating conditions; segregating the integrated semiconductor devices on the wafer by measuring a parameter of the devices; applying a stress test to the devices wherein test conditions of the stress test are adjusted based on the segregation parameter measurements of the devices; determining functionality of the devices at nominal operating conditions after the stress test; and, classifying the devices as failed if the devices do not function properly after the stress test.\nThe wafer includes devices having short and long gate channel widths, and the step of segregating the integrated semiconductor devices comprises determining which of the devices on the wafer have short gate channel widths. Applying a stress test pursuant to this method comprises applying a voltage higher than the device normal operating voltage to the devices. The step of determining which of the devices have short gate channel widths further comprises measuring the operational speed of the devices prior to applying the stress test.\nApplying a stress test comprises applying a first voltage at a value higher than the device normal operating voltage to the devices with a first measured operational speed, and a second voltage at a value lower than the first voltage to the devices with a second measured operational speed, the first operational speed being less than the second operational speed.\nMeasuring the operational speed of the devices may comprise using flush delay time measurements or a performance sort ring oscillator to determine the operational speed of the devices.\nThe wafer is fabricated to include n-type and p-type field effect transistors, and wherein the step of measuring the operational speed of the devices comprises the steps of:\ni) measuring drain-to-source current for the n-type and p-type field effect transistors;\nii) determining a drain-to-source current sum by summing the value of the n-type drain-to-source current with the absolute value of the p-type drain-to-source current;\niii) determining a device ISUM value by dividing the drain-to-source current sum by the gate channel width of the device; and,\niv) segregating the devices based on the ISUM value wherein the devices with lower ISUM values correspond to slower operational speeds and devices with higher ISUM values correspond to faster operational speeds.\nStep (iv), segregating the devices based on the ISUM value, further comprises assigning devices with ISUM values less than or equal to 700 xcexca/xcexcm as slower operational speed devices.\nThe current invention provides, in a third aspect, an apparatus for testing a wafer containing a plurality of integrated semiconductor devices including n-type and p-type field effect transistors with short and long gate channel widths, and electrical parameters comprising:\na computer containing therein a set of instructions adapted to perform the following functions:\na) measure the short and long gate channel widths of the semiconductor devices;\nb) segregate the semiconductor devices by the short and long gate channel width measurements;\nc) apply a voltage stress to the devices wherein the voltage stress is a function of the short and long gate channel width measurements; and,\nd) perform a functionality test to the devices after the voltage stress such that devices not passing the functionality test are labeled failed devices.\nIn a fourth aspect, the current invention provides a program storage device readable by machine, tangibly embodying a program of instructions executable by the machine to perform method steps for voltage screening semiconductor devices on a wafer, the method steps comprising:\na) determining functionality of the integrated semiconductor devices at nominal operating conditions;\nb) segregating the integrated semiconductor devices on the wafer by measuring a parameter of the devices;\nc) applying a stress test to the devices such that test conditions of the stress test are adjusted based on the segregation parameter measurements of the devices;\nd) determining functionality of the devices at nominal operating conditions after the stress test; and,\ne) classifying the devices as failed if the devices do not function properly after the stress test screen.\nLastly, in a fifth aspect, the current invention provides a method for testing integrated circuit semiconductor devices comprising the steps of:\na) providing a wafer containing a plurality of integrated semiconductor devices including n-type and p-type field effect transistors with short and long gate channel widths, and operational speed parameters;\nb) determining functionality of the integrated semiconductor devices at nominal operating conditions;\nc) segregating the integrated semiconductor devices on the wafer by measuring the operational speed parameters;\nd) applying a stress test to the devices wherein test conditions of the stress test are adjusted based on the operational speed measurements of the devices, such that a first voltage at a value higher than the device normal operating voltage is applied to the devices with a first measured operational speed, and a second voltage at a value lower than the first voltage is applied to the devices with a second measured operational speed, the first operational speed being less than the second operational speed;\ne) determining functionality of the devices at nominal operating conditions after the stress test; and,\nf) classifying the devices as failed if the devices do not function properly after the stress test."}
{"text": "The invention concerns a linear rolling bearing comprising a guide carriage with a U-shaped cross-section having a U-crossbar and two U-legs so that the guide carriage forms a carriage cavity and partially surrounds a guide rail while being slidably supported through balls on two longitudinal sides of the guide rail, each U-leg of the guide carriage having on an inner surface opposing the guide rail a ground raceway with an approximately quarter circle cross-section for the balls, the raceway being made by a grinding wheel whose diameter is larger than a diagonal dimension of the carriage cavity and whose axis of rotation is situated outside of the guide carriage and forms an acute angle with an axis of symmetry of the guide carriage.\nRaceways of carriages for profiled rail guides are made either using small vertical grinding wheels or several inclined grinding wheels arranged behind one another. Small grinding wheels result in unfavorable grinding values, short tool lives and poor surface quality. If several inclined wheels are used, the inner space of the guide carriage i.e., the carriage cavity likewise limits the size of the wheels. The advantage, however, is the inclined position of the grinding wheel. Therefore a limited but somewhat larger wheel is used.\nAt the same speed of rotation, a higher peripheral speed is obtained with a grinding wheel having a large diameter than with a wheel having a small diameter. With a large diameter, a higher grinding performance and a prolongation of the service life of the wheel are obtained. The rotational speed of the grinding spindle cannot be infinitely increased because this would lead to a destruction of its bearings.\nThe publication EP 0 318 980 B1 discloses a linear ball bearing having a guide carriage, designated as a bearing body, that has a leg section with an upper groove with an approximately semi-circular cross-section for load-bearing balls and a lower groove with an approximately quarter circle cross-section for load-bearing balls. A drawback of this lies in the continuation of the radius of the upper groove up to the vertical inner wall of the leg section. This necessarily means that for making the upper groove by grinding, a small grinding wheel must be used.\nThe publication DE 33 31 287 C2 discloses a method of grinding the bearing body of a linear ball bearing. In this method, two raceways for load-bearing balls are ground on the inner surface of a U-crossbar of a guide carriage using a grinding wheel with a large diameter which is placed inclined to the guide carriage. The return canals for the balls that are arranged in circuits are configured as bores in the guide carriage and have to be made in special work steps in addition to the raceways for the load-bearing balls.\nThe publication DE 33 04 895 C2 shows a linear ball bearing with bearing bodies in which each ball recirculating shoe or bearing body comprises attached sheet metal parts viz., a W-shaped ball retainer for the load-bearing balls and a cover with a U-shaped cross-section which outwardly closes ball-guiding grooves for the non-loaded returning balls. Due to these additional sheet metal parts, the structure and assembly of the ball recirculating shoe is complicated and expensive.\nThe publication DE 30 19 131 A1 discloses a linear ball bearing in which continuous longitudinal openings are provided in the load-bearing regions and the return regions for the balls in the guide elements that are made as retaining plates. Therefore each guide element not only has in its ball region a narrow outwardly oriented opening but is also open towards the machine element that acts as a guide carriage. The balls extend through these openings and are guided on longitudinal raceways that have to be made as grooves in the machine element so that additional machining and tool costs are incurred.\nIt is an object of the invention to provide a linear rolling bearing which can be made with low machining costs and a short machining time under good grinding conditions so that raceways with a very high surface quality can be realized.\nThis and other objects and advantages of the invention will become obvious from the following detailed description.\nThe invention achieves the above objects by the fact that a stop surface having a retaining contour for a guide member containing the balls is configured on a guide rail-distal outer surface of each U-leg of the guide carriage, and the raceway on one of the two U-legs and the stop surface having the retaining contour on the other of the two U-legs are made in common in one work step by the grinding wheel. Due to the simultaneous grinding of a raceway and a stop surface, the machining time of the guide carriage of the linear bearing is considerably shortened.\nFurther advantageous features of the invention and a detailed description of one example of embodiment of the invention shown in the drawings will be given hereinafter."}
{"text": "Digital Pathology, which can also be referred to as virtual microscopy or virtual pathology involves managing, analysing and interpreting digital information. In this field of technology the need to make inter-image comparisons and to train learning algorithms on large volumes of data represents a significant challenge. Whole slide microscope images may comprise large volumes of data. The applicants have recognised a need to automate upload and processing of such data without compromising the ability of such methods to provide objective comparisons between different tissue samples acquired using different systems,\nProcesses described herein may involve the generation of glass slides and converting these to digital pathology slides using digital pathology solutions. A digital slide scan is then generated which allows for high resolution viewing, interpretation and image analysis of digital pathology images. Growth in digital pathology solutions has totally transformed how research labs manage and interpret glass slides, with analysis and interpretation of histological images now conducted on a computer screen.\nGaining momentum globally, digital pathology is being used across health and pharmaceutical sectors, education and contract research organisations. With wide ranging applications the realised benefits of this sophisticated technology have encourage high growth in the market for digital pathology solutions, which by 2020 is estimated to be worth $5.7 billion.\nThe ability to provide accurate diagnosis is critical to the provision of healthcare. Biopsies to identify the presence of diseases such as cancer are a useful tool in such diagnosis. Accurate prognostic assessment is also critically important, because it enables action to be taken to counteract further development of disease. Microscope images of tissue samples have been used for these purposes for many years.\nLarge numbers of microscope imaging systems, each with their own particular characteristics, have been developed for this purpose. Whole slide imaging systems obtain digital images of entire microscope slides by scanning the field of view of a microscope across a macroscopic tissue sample to obtain a series of digital images. The resulting digital images can then be concatenated together to provide a single image, or image set, which describes the entire microscope slide. Partial images of slides can also be obtained by the same approach.\nPathologists involved in making diagnoses based on these kinds of images may rely on qualitative judgements. Such judgements may be based on their scientific knowledge and also on personal experience. This is necessarily a subjective process. As a result diagnoses, and prognostic assessments are not always reproducible—different pathologists may make different judgements based on identical images of tissue samples.\nIn making diagnostic judgements, the pathologist's task is made still more difficult because large tissue samples may be involved. In such cases many tens, or even hundreds of microscope images may need to be analysed from a single patient. This is particularly true where multiple tissue biopsies have been taken from a relatively large area of the body such as the prostate. These issues compound the problem of reproducibility because two different pathologists assessing the same patient's tissue sample may take into account features of different areas of different images of the same tissue sample.\nThe conditions under which a tissue sample was obtained, and the treatment of that sample before it was imaged (for example in terms of the concentration of stain applied to it), the imaging system used to acquire the image, and the presence of image artefacts may all cause variations between images. Although painstaking analysis is required, human pathologists are at least able intuitively to make allowances for such confounds. The subjective nature of assessment by human pathologists therefore, whilst problematic, at least provides one way to address these problems of inter-image variability. This need for intuitive judgement prevents straightforward automation of diagnostic and prognostic assessment of microscope images.\nThere are still further obstacles to overcome."}
{"text": "It has been known to measure voltages using optical means.\nA known technique for measuring high voltages is described in EP 0 682 261. The technique described in EP 0 682 261 relies on the electro-optical effect, in which an electrical field changes the refractive index or birefringence of a light-transmitting element, namely a bulk single crystal. The crystal is a Pockels-type electro-optical material that changes its refractive index or birefringence linearly with the applied electric field. The method described in EP 0 682 261 allows the voltage to be measured over the crystal. For this purpose, two orthogonally polarized light beams are passed through the crystal twice, and the returned light beams are brought to interference, e.g. in a polarizer or an interferometer, for measuring their mutual phase shift.\nOptical voltage sensors of this type for applications in air-insulated high-voltage substations commonly involve expensive dielectric insulation. The sensor of EP 0 682 261 makes use of the electro-optical effect in a bulk Bi4Ge3O12 (BOO) crystal. The voltage of a power line (up to several 100 kV) is applied to the crystal having a length typically between 100 mm and 250 mm. This results in very high electric field strengths near the crystal. To avoid dielectric breakdown, the crystal is mounted in a high-voltage insulator tube filled with SF6 gas under pressure.\nThe voltage sensor described in Ref. [2] as identified below uses several small electro-optical crystals which are mounted in a high-voltage insulator. The crystals measure the local electric fields. The sum of these local fields serves as an approximation of the line integral of the field. Here, the field strengths are lower and insulation with nitrogen at atmospheric pressure is sufficient. However, this technique involves such as extra measures permittivity-shielding to stabilize the electric field distribution and avoid excessive approximation errors. The present disclosure identifies particular references by numerical designation. The references are specifically identified herein following the detailed description of the exemplary embodiments."}
{"text": "1. Field of the Invention\nThis invention relates to a method, reagent, and reagent kit utilizing a surfactant-dye complex for the demonstration and determination of free fatty acids, compounds which hydrolyze to yield fatty acids, and hydrolase enzymes which act on substrates to produce fatty acids.\n2. Description of the Prior Art\nPrior methods for determining fatty acids in liquid samples include titration with alkali, determination of the copper salt of such acids after extraction with a suitable solvent, and chromatographic procedures. Such methods are time-consuming and, in general, are not adaptable to rapid, automated analysis techniques.\nFurther, various methods exist for demonstrating and determining compounds which are hydrolyzable to yield free fatty acids, including triglycerides, other cholesterol esters, and phospholipids. Prior methods for such determinations include enzymic hydrolysis methods, saponification methods, extraction methods, and chromatographic procedures.\nSuch techniques generally involve a number of processing steps which render the determinations time-consuming and costly, and/or require more than one reagent.\nHall U.S. Pat. No. 4,195,126 (Mar. 25, 1980), the disclosure of which is hereby incorporated by reference, discloses a method for determining free fatty acids and compounds which hydrolyze to yield free fatty acids.\nThe Hall patent describes an albumin-dye complex which is substantially free of endogenous free fatty acids and which has been pre-treated with a fatty acid. According to the method of the Hall patent, free fatty acids react with the albumin-dye complex to displace dye from the complex, yielding an albumin-fatty acid complex and free dye, the concentration of which is measurable by changes in absorbance or fluorescence. The change in absorbance is thus proportional to the concentration of fatty acid in the sample.\nThe Hall patent's complex may be used in determining a component in a sample which is capable of being hydrolyzed to yield fatty acid. The component is hydrolyzed by a hydrolase enzyme to yield a fatty acid, which is then reacted with the albumin-dye complex to liberate dye, the concentration of which can be directly measured by the change in absorbance or fluorescence.\nThe Hall patent also discloses that the complex can be used to determine the activity of a hydrolase enzyme in a liquid sample by adding to the sample a known amount of a compound which is capable of being hydrolyzed to yield a fatty acid, along with a reagent including the albumin-dye complex. The fatty acid thus liberated reacts with the complex to yield free dye, the concentration of which can be measured by observation of the absorbance change.\nThe albumin-dye complex of the Hall patent is limited in its practical utility, however, since one or more commercially useful hydrolase enzymes, such as Rhizopus arrhizus lipase, is conventionally provided in ammonium sulphate suspension. It has been found that ammonium sulfate interferes with the color development reaction of the Hall complex, thus necessitating desalting and lyophilizing of the lipase preparation obtained from the manufacturer. The lyophilized lipase is stable for only about 48 hours at 4.degree. C.\nFurther, it has been found that serum albumin depresses absorbance readings of the Hall reagent so that the absorbance of reagent-containing serum, without lipase, is lower than that of the reagent itself."}
{"text": "The present invention relates generally to medical devices and methods. More particularly, the invention relates to apparatus and methods for independently delivering a plurality of luminal prostheses within a body lumen.\nStenting has become an increasingly important treatment option for patients with coronary artery disease. Stenting involves the placement of a tubular prosthesis within a diseased coronary artery to expand the arterial lumen and maintain the patency of the artery. Early stent technology suffered from problems with restenosis, the tendency of the coronary artery to become re-occluded following stent placement. In recent years, however, improvements in stent design and the advent of drug-eluting stents have reduced restenosis rates dramatically. As a result, the number of stenting procedures being performed in the United States, Europe, and elsewhere has soared.\nStents are delivered to the coronary arteries using long, flexible vascular catheters, typically inserted through a femoral artery. For self-expanding stents, the stent is simply released from the delivery catheter, and it resiliently expands into engagement with the vessel wall. For balloon expandable stents, a balloon on the delivery catheter is expanded which expands and deforms the stent to the desired diameter, whereupon the balloon is deflated and removed.\nDespite many recent advances in stent delivery technology, a number of shortcomings still exist. For example, current stent delivery catheters are not capable of customizing the length of the stent in situ to match the size of the lesion to be treated. While lesion size may be measured prior to stenting using angiography or fluoroscopy, such measurements may be inexact. If a stent is introduced that is found to be of inappropriate size, the delivery catheter and stent must be removed from the patient and replaced with a different device of correct size. Moreover, current stent delivery devices cannot treat multiple lesions with a single catheter. If multiple lesions are to be treated, a new catheter and stent must be introduced for each lesion to be treated.\nAdditionally, currently available stent delivery devices are not well-adapted for treating vascular lesions that are very long and/or in curved regions of a vessel. Current stents have a discrete length that is relatively short due to their stiffness. If such stents were made longer, to treat longer lesions, they would not conform well to the curvature of vessels or to the movement of vessels on the surface of the beating heart. On the other hand, any attempt to place multiple stents end-to-end in longer lesions is hampered by the inability to maintain appropriate inter-stent spacing and to prevent overlap of adjacent stents. Such shortcomings in the prior art are addressed by the inventions described in U.S. patent application Ser. No. 10/412,714, which is hereby fully incorporated by reference, and Ser. No. 10/637,713, which was previously incorporated by reference.\nEven with improvements such as those described in the above-referenced patent applications, further improvements in stent delivery devices and methods are still being sought. For example, before a coronary stent is deployed in a stenotic lesion, the physician will typically first dilate the lesion with an angioplasty balloon. Following such “predilatation,” the angioplasty catheter is removed and a stent delivery catheter is advanced to the treatment site to deploy the stent. One of the significant advantages of the stent delivery systems described in U.S. patent application Ser. Nos. 10/412,714 and 10/637,713, incorporated above, is the ability to treat multiple lesions at different locations without removing or replacing the catheter. Such a stent deliver system may be positioned at a first lesion for deployment of a first stent of a desired length, then moved to a second site where a second stent of a different length may be deployed. This may be repeated for multiple lesions without exchanging catheters, which saves time and eliminates the inefficiency of using multiple catheters. Such efficiencies are reduced, however, if it is necessary to use an angioplasty catheter to predilate lesions and a separate stent delivery catheter to deliver stents. If separate predilatation and stent delivery catheters are used, it may often be necessary to exchange, or “swap out,” the two catheters multiple times during a stenting procedure.\nTherefore, it would be desirable to have stent delivery systems that could be used to predilatate lesions without requiring a separate angioplasty catheter. Ideally, such stent delivery systems would allow for separate predilatation of multiple and/or long lesions as well as separate stent deployment at those lesions, without requiring any catheters to be exchanged. Preferably, such systems would also enable a user to adjust the length of the predilatation device to match the length of the lesion to be treated. At least some of these objectives will be met by the present invention."}
{"text": "A file server is a computer that provides file service relating to the organization of information on storage devices, such as disks. The file server or filer may be embodied as a storage system including a storage operating system that implements a file system to logically organize the information as a hierarchical structure of directories and files on the disks. Each “on-disk” file may be implemented as a set of disk blocks configured to store information, such as text, whereas each directory may be implemented as a specially-formatted file in which information about other files and directories are stored.\nAs used herein, the term storage operating system generally refers to the computer-executable code operable on a storage system that manages data access and client access requests and may implement file system semantics in implementations involving filers. In this sense, the Data ONTAP™ storage operating system, available from Network Appliance, Inc. of Sunnyvale, Calif., which implements a Write Anywhere File Layout (WAFL™) file system, is an example of such a storage operating system implemented as a microkernel within an overall protocol stack and associated disk storage. The storage operating system can also be implemented as an application program operating over a general-purpose operating system, such as UNIX® or Windows®, or as a general-purpose operating system with configurable functionality, which is configured for storage applications or debugging/servicing applications as described herein.\nA filer's disk storage is typically implemented as one or more storage volumes that comprise physical storage disks, defining an overall logical arrangement of storage space. Available filer implementations can serve a large number of discrete volumes (150 or more, for example). A storage volume is “loaded” in the filer by copying the logical organization of the volume's files, data and directories into the filer's memory. Once a volume has been loaded in memory, the volume may be “mounted” by one or more users, applications, devices, etc. permitted to access its contents and navigate its namespace. As used herein, a volume is said to be “in use” when it is loaded in a filer's memory and at least one user, application, etc. has mounted the volume and modified its contents.\nA filer may be configured to operate according to a client/server model of information delivery to thereby allow many clients to access files stored on a server, e.g., the filer. In this model, the client may comprise an application, such as a file-system protocol, executing on a computer that “connects” to the filer over a computer network, such as a point-to-point link, shared local area network (LAN), wide area network (WAN), or virtual private network (VPN) implemented over a public network such as the Internet. Communications between the filer and its clients are typically embodied as packets sent over the computer network. Each client may request the services of the filer by issuing file-system protocol messages formatted in accordance with a conventional file-system protocol, such as the Common Internet File System (CIFS) or Network File System (NFS) protocol.\nWhen the filer runs the storage operating system or other applications, a problem/error in the execution of programming code of the storage operating system or the other applications may occur at any given time. Of particular interest are three types of problem situations that may occur while executing applications on the filer: 1) when an unrecoverable exception occurs causing a “crash” and reboot of the filer, 2) when the performance of the filer is significantly reduced, and 3) when the filer is “wedged” and can not receive and perform administrative commands. In the first situation, an exception in the operating system has occurred that renders the filer inoperative and the filer must be rebooted (the filer restarted and the operating system reloaded). In the second situation, the filer read/write operations for client requests have significantly slowed due to some operating system error (e.g., a memory leak where the filer's memory resources are being allocated but are not being released after usage). In the third situation, the filer is “wedged” when it may or may not be able to perform read/write operations, but can not receive and perform administrative commands (command line instructions), i.e., is apparently unresponsive to all or certain types of commands or operations, in particular administrative commands.\nIn all three problem situations, a reboot of the filer can be performed. In the first situation, a reboot is automatically performed. In the second and third situations, a manual reboot can be performed (e.g., by issuing an administrative command or pressing the filer's reset button). During the reboot, the filer typically performs a reboot/shut-down procedure that includes a corefile routine that generates a corefile (core dump) that is stored to the filer's memory. The corefile comprises a static image of the memory content/data and state of the filer at the time the corefile routine is performed. The corefile can then be analyzed by a debugging program (debugger) operated by a programmer to determine the problem/error that occurred during the execution of the operating system or other application and to help develop programming code that will avoid the problem/error in the future.\nThe corefile routine creates a corefile comprising a corefile header and data that is copied from the filer's memory (referred to herein as “filer memory data”). The corefile header comprises corefile metadata (data describing the corefile and the filer) and a set of memory range descriptors that provide an address mapping table between the filer's memory addresses and the corefile addresses. Typically, only particular areas of the filer memory data are copied to the corefile. These areas are generally those areas of filer memory data that are accessible and important for debugging purposes.\nTypically, the debugger resides on a client administering computer that is remote from the filer and receives the corefile through a network that connects the administering computer and filer. However, corefiles can be very large in size (due to the copied filer memory data) and require long upload times to the administering computer. As such, a Core Daemon Protocol (CDP) can be used to allow remote analysis of the corefile by the debugger without requiring uploading of the corefile to the administering computer where the debugger resides.\nThe Core Daemon Protocol is a simple file-access protocol specific for corefiles that allows for the remote reading and retrieving of parts of the corefile. As known in the art, the Core Daemon Protocol specifies a set of rules (e.g., data packet format, sequence of events to occur, etc.) for communication between the administering computer (that runs the debugger) and the filer for remotely accessing the corefile stored on the filer. The Core Daemon Protocol provides simple operations such as Open File Operation (Open(“corefile name”)), Read L bytes at address/offset A (Read(Offset, Length)), and Close File Operation (Close) to allow remote open, read, and close file operations to a specified file (“corefile name”). The corefile routing typically stores the corefile to a predetermined path/location on the filer (e.g., the filer memory's root directory) with a predetermined filename. The debugger/programmer will have knowledge of the predetermined path/location and filename so the debugger/programmer can locate and access the corefile.\nA core daemon (remote core analysis daemon) is a program used to implement the Core Daemon Protocol. The core daemon program works in conjunction with the debugger and accesses the corefile and responds to requests from the debugger in accordance with the Core Daemon Protocol. In particular, the debugger submits requests for data in the corefile to the core daemon program which retrieves the data from the corefile and sends the data to the debugger. As such, the debugger can receive and analyze data of the corefile without requiring uploading of the entire corefile to the remote administering computer. The core daemon program typically resides on a support console computer connected to the filer. The support console typically runs only specialized administrative applications (such as the core daemon program) for administering and servicing the filer. The support console is configured to access the filer's file system and files stored on the filer.\nThe two above-mentioned methods for debugging the filer's operating system or other applications (i.e., uploading the corefile to the administering computer or performing remote corefile analysis) are considered offline or non-live debugging since the generated corefile is a static image of the filer's memory at a previous time. Another method for debugging the filer is debugging the current condition of the file as it is running (i.e., live debugging). Typically, in live debugging, a “crash” type error has not occurred, but rather, the filer's read/write operations have slowed significantly or the filer has “wedged.” In live debugging, a debugger is operated by a programmer on a remote administering computer. The debugger gains control of the filer's operating system and “lock steps” the operations of the filer (i.e., does not allow the filer to perform other operations while waiting for commands from the programmer). The filer is unlocked to perform other operations only when the programmer allows the filer to do so. In this way, while the programmer may be pondering a data result, and does not unlock the filer, the filer can not perform other operations.\nThe current methods for debugging the filer has several disadvantages. The non-live debugging methods that require a corefile to be generated are disruptive in that they require a reboot procedure that disables all filer processes and shuts down the filer for a particular down-time period (until the filer is restarted and able to again perform operations). In addition, the creation of the corefile can significantly increase the down-time period as the corefile routine copies and stores large amounts of filer memory data to the corefile. Although there are several methods for reducing the size of the corefile, the corefile is still relatively large in size (e.g., typically around 6-12 GB for a 30 GB filer memory). Live debugging is also disruptive in that the operations of the filer are shut down by the debugger while waiting for commands from the programmer. Only when the filer is unlocked by the programmer can the filer perform read/write operations. As such, there is a need for a less-disruptive method for debugging a computer system, such as a filer."}
{"text": "1. Field of the Invention\nThe present invention relates to redundant power distribution systems, and more particularly to replacing a power ORing diode with a FET and a controller that operates the FET as a pass device.\n2. Description of the Related Art\nIn a multiple supply fault tolerant redundant power distribution system, the paralleled power supplies are operated to contribute equally to the load current through various power/current sharing schemes. Regardless of the particular current sharing scheme, a common design practice is to include discrete ORing power diodes to protect against reverse current flow should one of the power supplies develop a catastrophic failure resulting in a high current path to ground. In addition, reverse current can occur if the current sharing scheme fails and an individual power supply voltage falls significantly below the others.\nAlthough discrete ORing diodes have been used for some time and are relatively inexpensive to implement, there are several drawbacks to their use. One problem is the associated increasing power dissipation loss in the ORing diodes as power requirements for the system increases. The power loss across a typical ORing diode at 20 Amperes (A) is about 14 Watts (W), which becomes significant when multiplied by the number of redundant power supplies.\nIt is desired to provide similar functionality as the ORing diode in a redundant power distribution system without the concomitant loss of power at high loads."}
{"text": "1. Field of the Invention\nThe present invention relates to a display unit in which light emitting pixels are arrayed in a matrix form, and an image signal having a half tone can be displayed.\n2. Description of the Prior Art\nFIG. 1 is a block diagram showing a conventional display unit disclosed in Japanese Patent Publication (Kokoku) No. 2-709. In FIG. 1, reference numeral 1 denotes a matrix display panel such as a fluorescent character display tube, 2 is a driver to drive row electrodes of the matrix display panel 1, and 3 is another driver to drive column electrodes of the matrix display panel 1. Display elements (light emitting elements) are disposed at intersections of the row electrodes and the column electrodes, and are turned ON by driving the corresponding row electrode and the corresponding column electrode. Reference numeral 4 denotes a shift register to place an ON/OFF signal from a changing portion 15 on the corresponding column electrode, 5 is a row electrode control circuit to create a drive signal for the row electrode, and 6 is a column electrode control circuit to create a drive signal for the column electrode. The column electrode control circuit 6 latches output of the shift register 4 for a predetermined time interval depending upon a latch signal from a timing generating circuit 8.\nReference numeral 9 denotes a display control unit to write display data onto a memory 7, 10 is a selector to feed the memory 7 with any one of a read address and a write address, 11 is a read/write control circuit to feed the selector 10 with a switching signal, 12 is a clock generating circuit to feed a clock signal to the timing generating circuit 8 and the read/write control circuit 11, 13 is a read address counter to generate the read address, and 15 is the changing portion to transform data read from the memory 7 into the ON/OFF signal for pixel.\nThe matrix display panel 1 is provided with a structure as shown in FIG. 2. That is, the display elements are disposed at intersections of a group of signal lines for feeding signals to the column electrodes X.sub.1 to X.sub.m and a group of signal lines for feeding signals to the row electrodes Y.sub.1 to Y.sub.m. Thus, the display elements are controlled according to a combination of the signals fed to the two groups of signal lines.\nA description will now be given of the operation. The display control unit 9 outputs the display data, an address corresponding to a position at which the display data is displayed, and a timing signal in synchronization with the display data. The timing signal is inputted into the read/write control circuit 11. The read/write control circuit 11 receives the timing signal as input, and thereafter outputs a signal by which the selector 10 is switched over to the write address. Therefore, the address from the display control unit 9 is fed to the memory 7 as the write address. As a result, the display data can be stored in the memory 7 at an area specified by the write address.\nFurther, the display data stored in the memory 7 is read out according to the read address which is produced in the read address counter 13. The read address counter 13 generates an X address (a column address) fed to the column electrode control circuit 6, a Y address (a row address) fed to the row electrode control circuit 5, and a comparison signal B. The X address and the Y address are fed into the memory 7 through the selector 10. In this case, the selector 10 selects the read address according to control by the read/write control circuit 11. Consequently, it is possible to output data from the memory 7 according to the read address including the X address and the Y address.\nFurther, the Y address is sent to the row electrode control circuit 5. The row electrode control circuit 5 decodes the Y address so as to drive the row electrodes Y.sub.1 to Y.sub.m of the matrix display panel 1 through the driver 2. The data read from the memory 7 is sent to the changing portion 15. The changing portion 15 compares the data with the comparison signal B produced in the read address counter 13 so as to generate an ON/OFF signal according to the result of comparison. When the data is greater than the signal B, the ON signal is generated. When the data is equal to or less than the signal B, the OFF signal is generated. Here, `1` shows the ON signal and `0` shows the OFF signal. The ON/OFF signal is disposed in the shift register 4 according to a shift clock from the timing generating circuit 8. Further, the signal is latched by the column electrode control circuit 6 in response to a latch signal from the timing generating circuit 8. In response to the latched signal, the column electrode control circuit 6 drives the column electrodes X.sub.1 to X.sub.m through the driver 3.\nIn such a display unit, an image is displayed by sequentially and periodically driving the row electrodes Y.sub.1 to Y.sub.m, and by switching operation of a signal which is fed to the column electrodes X.sub.1 to X.sub.m in synchronization with the driving of the row electrodes Y.sub.1 to Y.sub.m. The display elements can exclusively provide binary representation, that is, any one of ON and OFF. However, when an image signal having a half tone should be displayed, the data in the memory is read out a predetermined number of times, and a cumulative elapsed ON time in each display element is controlled, thereby realizing gray image display.\nA more detailed description will now be given of a display operation by way of, as an example, the matrix display panel 1 including a device having a 4-by-4 pixel array. Pixel positions in the matrix display panel 1 correspond to memory addresses in a one-to-one manner, resulting in a relationship as shown in FIGS. 3 and 4. FIG. 3 shows the pixel positions, and FIG. 4 shows the memory addresses corresponding to the respective pixel positions. FIG. 5 shows illustrative display data stored in the memory 7 at the respective addresses. As shown in FIG. 6, the addresses of the memory 7 can be classified into the X address and the Y address. FIG. 9 is a timing diagram showing row electrode driving timing and column electrode driving timing. FIG. 10 is an explanatory view showing the result of the driving.\nWhen the display data corresponding to each pixel includes a 4-bit structure, it is possible to display a 16-scale gray image (because of 2.sup.4 =16). In this case, as shown in FIG. 7, the read address counter 13 includes a 2-bit counter 13a for generating the X address, a 4-bit counter 13b for generating the comparison signal B, and a 2-bit counter 13c for generating the Y address. The 2-bit counter 13b counts the clock signal from the read/write control circuit 11 to sequentially output values of 0 to 3 corresponding to the column electrodes X.sub.1 to X.sub.4. The 4-bit counter 13b counts a carry signal of the 2-bit counter 13a. In the 4-bit counter 13b, an initial value is set to zero. The 2-bit counter 13c counts a carry signal of the 4-bit counter 13b. The initial value of zero in the 2-bit counter 13c corresponds to the row electrode Y.sub.1.\nAccordingly, with the comparison signal B of zero and the Y address of zero, four data stored in the first row shown in FIG. 5 are sequentially read out. As shown in FIG. 8, the changing portion 15 comprises, for example, a comparator 15a. The comparator 15a compares display data A with the comparison signal B so as to output an ON signal when A is greater than B, or output an OFF signal when A is equal to or less than B. As shown in FIGS. 5 and 6, the first row sequentially contains data of \"15 (in decimal notation),\" \"10 (in decimal notation),\" \"12 (in decimal notation),\" and \"0\". Since the comparison signal B is zero, the changing portion 15 sequentially outputs \"1,\" \"1,\" \"1,\" and \"0\". Signals from the changing portion 15 are sequentially inputted into the shift register 4.\nWhen the respective signals in the first row are inputted into the shift register 4, the timing generating circuit 8 outputs the latch signal such that the contents of the shift register 4 can be latched by the column electrode control circuit 6. The column electrode control circuit 6 drives the column electrodes X.sub.1 to X.sub.4 through the driver 3 according to the contents of the shift register 4. At the time, a value of zero is fed into the row electrode control circuit 5 as the Y address so that the row electrode control circuit 5 can drive the row electrode Y.sub.1 in the first row through the driver 2. Thus, as shown by a period t1 in FIGS. 9 and 10, display elements at intersections of the first row and the first to third columns are turned ON, and a display element at an intersection of the first row and the fourth column is not turned ON. In FIG. 10, a hollow ring denotes a display element which is turned ON, and a numeral in the hollow ring denotes an ON time interval. Further, reference numerals t1 to t64 correspond to reference numerals t1 to t64 in FIG. 9.\nNext, by the read/write control circuit 11, a count value of the 2-bit counter 13a is reset to zero, and a count value of the 4-bit counter 13b is set to one. A count value of the 2-bit counter 13c is zero. Consequently, the four display data in the first row are read from the memory 7 again, and the display data are compared with the comparison signal B (=1) in the changing portion 15. According to the result of comparison, the changing portion 15, the shift register 4, the column electrode control circuit 6, and row electrode control circuit 5 are operated as set forth above. Therefore, as shown by a period t2 in FIGS. 9 and 10, the display elements are turned ON at the intersections of the first row and the first to third columns, and the display element at the intersection of the first row and the fourth column is not turned ON.\nWhen the count value of the 4-bit counter 13b reaches 15 (in decimal notation) after count from 0 to 3 is repeated in the 2-bit counter 13a, the changing portion 15 compares the data in the first row from the memory 7 with the comparison signal B (=15). The first row sequentially contains the data of \"15 (in decimal notation),\" \"10 (in decimal notation),\" \"12 (in decimal notation),\" and \"0.\" Because of the comparison signal of 15, the changing portion 15 sequentially outputs \"0,\" \"0,\" \"0,\" and \"0.\" As a result, as shown by a period t16 in FIGS. 9 and 10, all the display elements at the intersections of the first row and the first to fourth columns are not turned ON. As set forth above, comparison is made the number of times identical with the number of gray scale for each row data. As shown at the leftmost end of the first row in FIGS. 9 and 10, according to the result of comparison, the display element corresponding to the address \"0\" is held ON for a period of 15 unit times, the display element corresponding to the address \"1\" is held ON for a period of 10 unit times, and the display element corresponding to the address \"2\" is held ON for a period of 12 unit times. In such a manner, a gray image can visually be displayed according to the ON time periods.\nNext, the count values of the 2-bit counter 13a and the 4-bit counter 13b are reset to zero, thereby setting the count value of the 2-bit counter 13c to one. Thus, the above processing is performed for data in the second row corresponding to the Y address \"1\". As a result, as shown at the leftmost end of the second row in FIGS. 9 and 10, the display element corresponding to the address \"4\" is held ON for a period of 8 unit times, the display element corresponding to the address \"5\" is held ON for a period of 4 unit times, the display element corresponding to the address \"6\" is held ON for a period of 2 unit times, and the display element corresponding to the address \"7\" is held ON for a period of 14 unit times.\nThe above processing is similarly carried out for data in the third row (the Y address of 2) and the fourth row (the Y address of 3), resulting in completion of the display operation in one frame. In the conventional display unit shown in FIG. 1, all the ON signals have a constant time interval as shown in FIG. 11, and a pulse number control method is employed in which the gray image display is realized according to the number of the ON signals.\nOn the other hand, another display unit is known in which pulse width control is made. In such a display unit, as shown in FIG. 12, a changing portion 15 has a selector 15b which can select one bit among 4-bit display data from a memory 7 according to a select signal S. The selection of the selector 15b results in an ON/OFF signal. Further, a timing generating circuit 8 is controlled to vary a time interval of the ON signal according to a bit weight (selected from weights 1, 2, 4, and 8). FIG. 13 is a timing diagram showing row electrode driving timing and column electrode driving timing. FIG. 14 is an explanatory view showing the result of driving.\nIn this case, a counter for generating an X address in a read address counter 13 sequentially feeds the memory 7 with values of 0 to 3. Output timing of the values of 0, 1, 2, and 3 do not have a constant interval, but have an interval according to the bit weight. A counter for generating the select signal S is positioned at a subsequent stage of the counter for generating the X address, and is incremented by one after the counter for generating the X address outputs all the values of 0 to 3. Further, another counter for generating a Y address is positioned at a subsequent stage of the counter for generating the select signal S, and is incremented by one after the counter for generating the select signal S outputs all values of 0 to 3.\nTherefore, four data in the first row in the memory 7 are respectively read out four times (according to timing of the periods t1 to t4), and are compared with the select signal S for each time. The first select signal S is set to zero, the second select signal S is set to one, the third select signal S is set to two, and the fourth select signal S is set to three. In the first comparison, the changing portion 15 selects and outputs the least significant bit among the 4-bit data from the memory 7. In the second comparison, the changing portion 15 selects and outputs the first bit among the 4-bit data from the memory 7. In the third comparison, the changing portion 15 selects and outputs the second bit among the 4-bit data from the memory 7. Finally, in the fourth comparison, the changing portion 15 selects and outputs the most significant bit among the 4-bit data from the memory 7.\nIn such a manner, when, for example, data at the address \"0\" shown in FIG. 6 is outputted from the memory 7 four times, the changing portion 15 outputs \"1,\" \"1,\" \"1,\" and \"1\" according to the timing of the periods t1 to t4. Since the output respectively have weights 1, 2, 4, and 8, in display modes t1 to t4 in the first row shown in FIG. 14, display is made for time intervals obtained by the leftmost values corresponding to the address \"0\". When data at the address \"1\" is outputted from the memory 7 four times, the changing portion 15 outputs \"0,\" \"1,\" \"0,\" and \"1\" according to the timing of the periods t1 to t4. Thus, in the display modes t1 to t4 in the first row shown in FIG. 14, display is made for time intervals obtained by the second values from the left corresponding to the address \"1.\"\nFor data at the address \"2\" and at the address \"3,\" processing identical with the above processing is carried out by the changing portion 15, the timing generating circuit 8, and so forth. In such a manner, signals corresponding to signals in the periods t1 to t4 in FIG. 13 are fed to a row electrode Y.sub.1 and column electrodes X.sub.1 to X.sub.4. As shown by the leftmost display mode in the first row in FIG. 14, a display element corresponding to the address \"0\" is held ON for a period of 15 unit times, a display element corresponding to the address \"1\" is held ON for a period of 10 unit times, and a display element corresponding to the address \"2\" is held ON for a period of 12 unit times.\nSubsequently, the Y address fed to both the memory 7 and the row electrode control circuit 5 is set to one, and the same processing as the above processing is carried out for data at the addresses \"4\" to \"7.\" Thus, in case of data as shown in FIG. 6, signals corresponding to signals in periods t5 to t8 in FIG. 13 are fed to a row electrode Y.sub.2 and column electrodes X.sub.1 to X.sub.4. The same operation is carried out with respect to Y addresses \"2\" and \"3.\"\nIn the display unit, the ON signals are set to have four types of time intervals, and a gray image can be realized by a combination of the ON signals as shown in FIG. 15. For example, in case of display data of weight 8, an ON signal having a time interval of 8 is fed to the column electrodes X.sub.1 to X.sub.4. Alternatively, in case of display data of 7, an ON signal having a time interval of 1, an ON signal having a time interval of 2, and an ON signal having a time interval of 4 are respectively fed to the column electrodes X.sub.1 to X.sub.4.\nIn the display unit employing the pulse number control method, the display data must be read out from the memory 7 the number of times identical with the number of gray scale before one frame can be displayed. Hence, when an increase in the number of bits in the display data increases the number of gray scale, it is necessary to operate the memory 7 and the control circuit at a higher velocity. As a result, the increase in the number of gray scale is limited.\nIn the display unit employing the pulse width control method, the number of times the display data is read from the memory 7, required to display the one frame, becomes equal to the number of bits in the display data. Thus, the pulse width control method requires considerably less number of times of read than the pulse number control method. However, for a short time interval such as the period t1 or t5, it is necessary to read out the display data expressed as the number of bits per data by the number of pixels per row (for example, four bits by four pixels in case of a 16-scale image including a 4-by-4 pixel array). Consequently, it is also necessary to operate the memory 7 and the control circuit at a higher velocity, and the increase in the number of gray scale is limited.\nFurther, due to influences such as response velocities of drivers 2 and 3 and light emitting elements, or ripple noise in power supply voltage, it becomes difficult to, for example, realize a pulse width corresponding to a time interval of 7 by summing up the ON signals respectively having the time interval of 1, the time interval of 2, and the time interval of 4. As the number of scale is further increased, such a tendency becomes more pronunced. When the scale number is increased, a value of the display data can not easily be made proportional to brightness of the display panel because of the following reason. That is, when the number of gray scale is increased to reduce an interval of the ON signal, sometimes supply of the ON signal fails to turn ON the display element by a response rise delay time and so forth due to the influences such as response velocity of the display element. That is, at times, the value of the display data can not easily be made proportional to brightness of display.\nIn addition, both the control circuit in the pulse number control method and the control circuit in the pulse width control method feed the ON signals to the column electrodes and the row electrodes according to the same timing. As a result, an expensive circuit having a higher response velocity is required for both a column electrode drive circuit and a row electrode drive circuit."}
{"text": "This invention relates to key loops such as those used in releasably supporting keys in depending relation from a key head affixed within various types of key cases. The type of key case typical of those to which the key loops of the present invention may be adapted is shown in U.S. Pat. Nos. 4,058,602 and 4,086,796.\nKey support loops generally comprise a spring metal rod or wire stock formed into a loop with an access opening at the top. The loop may be a \"tear drop\" or \"pear\" shape with an enlarged head or ball for attachment to a key head or bar mounted across the inside top of the key case. The free end of the resilient loop terminates adjacent the shank portion thus forming an access opening or mouth adapted to receive a key which may be fitted onto the loop by inserting the apertured head of the key onto the free upper end of the loop and pressing downwardly to wedge open the mouth portion of the loop.\nThose who are familiar with these types of key cases are aware that attachment of a key to a support loop usually requires a fair degree of care and dexterity. Even after careful attachment of the keys, however, it is not uncommon to find that one or more keys has become detached from such conventional key loops. This happens by the key head working its way toward the mouth of the key loop and wedging apart the mouth portion of the loop so that the key may become separated from the key case.\nPrior art attempts to render key loops more secure from such key loop separation have involved the use of closure attachments for spanning the key receiving mouth portion of the loop. Another approach has been to provide key loops of relatively complex configuration in which the upper ends of the loop are of split construction and the loop is opened and closed by twisting its upper ends. It will be realized that the manufacture of such loop structures involve greater manufacturing time than a simple loop and are thus more costly. In addition, such loops must be removed from the key loop support bar in order to attach or remove a key therefrom.\nIt is a principal object of the present invention to provide a key hook which overcomes the deficiencies of the prior art.\nAnother object is to provide a key hook wherein the problem of accidental key detachment is essentially eliminated.\nYet another object of the present invention is to provide a key hook on which a key can be attached easily and conveniently.\nIt is a further object of the present invention to provide a key hook which is simple in construction and economical to manufacture."}
{"text": "1. Field of the Invention\nThe present invention relates to optical glass, and more particularly, to optical glass having low melting point and which is suitable for a press lens.\n2. Description of the Related Art\nThere has been, as shown in Japanese Patent Application Laid-open Specification Nos. 56-59641, 56-149343 and the like, a study for a press lens which can be produced by press molding a glass heated above a softening point in a metal mold without having to subsequently grind and polish glass.\nWhen producing such kind of press lens, there has been problem regarding the melting point of optical glass employed. When the melting point is high, the metal mold used for molding can readily cause roughness of the surface by oxidation and the like and the life of the metal mold is thus shortened. Accordingly, it may be said that such optical glass having a low melting point is advantageous.\nOn the other hand, such press lenses are generally applied to minute parts so that the specific gravities of these presses must be as small as possible. It is preferable to decrease the thermal expansion coefficient since the contraction amount of the press lens must be small. In addition, the press lens naturally needs to have excellent weather-resistent characteristics.\nAs to optical glass, lenses having an intermediate index of refraction of 1.65 to 1.75 are widely available commercially. There has been a need for producing a lens having an intermediate index of refraction by press molding. In view of foregoing problems, when the conventional glass compositions are examined, SK glass and the like are difficult for press lenses since they have a higher melting point than that of SF series optical glass.\nIn conventional SF series glass, there have been known glass having components of SiO.sub.2, Al.sub.2 O.sub.3, Na.sub.2 O, K.sub.2 O, TiO.sub.2, PbO and As.sub.2 O.sub.3. This glass has an intermediate index of refraction of 1.75 to 1.9, specific gravity of not more than 5 g/cm.sup.3, and excellent weather-proofness. In such conventional glass, the life of the metal model can not be lengthened since many of these glasses have high melting points of not less than 440.degree. C.\nIn addition, in the conventional SF series, there has been no glass having an index of refraction approximately in the range of 1.65 to 1.75."}
{"text": "Today, new types of communication networks are being developed for the transfer of information in circuit-switched channels of time division multiplexed networks. In such a network, a recurrent frame of the network is divided into time slots, and said circuit-switched channels are established to comprise respective one or more time slots of said recurrent frame. Hence, each channel will be allocated a respective set of one or more time slots within said recurring frame, thereby gaining exclusive write access to said respective set of time slots.\nOne example of such a network is the so-called DTM network (DTM—Dynamic Synchronous Transfer Mode). DTM is a broadband network architecture (see e.g. Christer Bohm, Per Lindgren, Lars Ramfelt, and Peter Sjödin, The DTM Gigabit Network, Journal of High Speed Networks, 3(2), 109-126, 1994, and Lars Gauffin, Lars Hak{dot over (a)}nsson, and Björn Pehrson, Multi-gigabit networking based on DTM, Computer Networks and ISDN Systems, 24(2), 119-139, April 1992).\nPrior art describes rules for deciding, at a node of such a network, when the node shall give away write access to time slots to other nodes, i.e. to deallocate time slots in the favor of other nodes. For example, as disclosed in WO 97/36402, a node will give away time slots if it has any free slots available, i.e. slots not allocated to any one of the end users served by said node. However, in many cases there will be a need for more flexible rules, providing more options in the control of allocation and deallocation of time slots. Also, in this prior art, a resource shortage problem will occur if there are no free slots available in the frame of interest to respond to a request for additional slots transfer capacity."}
{"text": "1. Field of the Invention\nThis invention relates to a process for reforming a naphtha stream over a bifunctional, metal loaded aluminosilicate catalyst.\n2. Description of Related Art\nNaphtha streams emerging from some petrochemical refining processes generally comprise a mixture of C.sub.5 to C.sub.13 hydrocarbons which include about 15-40 wt % of C.sub.6 to C.sub.11 aromatic compounds and the balance mostly a mixture of C.sub.5 to C.sub.11 aliphatic hydrocarbons, including mixed paraffins and mixed olefins. The bulk of the naphtha stream, e.g. at least about 25 wt %, comprises C.sub.5 to C.sub.9 hydrocarbons, many of which boil near the boiling range of the benzene/toluene/xylene (BTX) fractions present in the naphtha. This makes it difficult to extract the more valuable BTX components from the naphtha by conventional distillation techniques. Therefore, a solvent extract process is required to purify the BTX component for chemical uses, which adds to the cost of recovery of BTX components. The naphtha stream from the refinery also contains sulfur contaminants such as elemental sulfur, alkyl sulfides and sulfur compounds such as benzothiophenes.\nA conventional procedure for both upgrading the naphtha and removal of the sulfur is to subject the naphtha stream to hydrodesulfurization (HDS) wherein the stream is contacted at high temperatures and in the presence of hydrogen with a desulfurization catalyst such as a sulfided cobalt or nickel/molybdenum catalyst. In addition to sulfur removal, the HDS process results in some aromatization and cracking of the C.sub.5 -C.sub.9 hydrocarbons present in the naphtha, thereby facilitating the ability to separate BTX components form the hydrorefined product. However, the HDS process consumes large quantities of hydrogen, e.g. up to 10,000 SCF/B, rendering it an expensive process. Nonetheless, it is desirable to remove sulfur because it tends to poison conventional catalysts which are used to reform naphtha, e.g. noble metal loaded aluminate or aluminosilicate catalysts.\nIt is also known in the art that HDS treated naphtha streams, either prior to or after removal of substantial quantities of the BTX fraction, can be subjected to catalytic reforming to further enhance the aromatics content of the naphtha. In a typical reforming process, the reactions include dehydrogenation, dehydrocyclization, isomerization, and hydrocracking. The dehydrogenation reactions typically include dehydroisomerization of alkylcyclopentanes to aromatics, the dehydrogenation of paraffins to olefins, the dehydrogenation of cyclohexanes to aromatics and the dehydrocyclization of acyclic paraffins and acyclic olefins to aromatics. The aromatization of the n-paraffins to aromatics is generally considered to be the most important because of the high octane rating of the resulting aromatic product. The isomerization reactions included isomerization of n-paraffins to isoparaffins, the hydroisomerization of olefins to isoparaffins, and the isomerization of substituted aromatics. The hydrocracking reactions include the hydrocracking of paraffins and hydrodesulfurization of any sulfur compounds remaining in the feed stock.\nIt is well known that several catalysts are capable of reforming petroleum naphthas and hydrocarbons that boil in the gasoline boiling range. Examples of known catalysts useful for reforming include platinum (optionally with the addition of rhenium or iridium) on an alumina support, platinum on zeolites of small pore size such as type X and Y (provided the reactants and products are sufficiently small to flow through the pores of the zeolites), and platinum on zeolite KL supports as disclosed in U.S. Pat. No. 4,987,109 and WO91/06616. Catalytic reforming of essentially sulfur-free aliphatic hydrocarbons using a zinc or gallium loaded ZSM-5 catalyst is also disclosed in U.S. Pat. Nos. 3,756,942, 4,180,689, 4,490,569 and 4,933,310, as well as by Fukase et al., \"Catalysts in Petroleum Refining And Petrochemical Industries 1995,\" 1996, pp. 456-464. Catalytic reforming of a sulfur-free naphtha using a ZMS-5 catalyst loaded with both platinum and another metal such as zinc or gallium is disclosed in WO96/03209.\nHowever, none of these references teach reforming of a naphtha stream which has not been desulfurized and which contains high concentrations of C9+ aromatics above about 5 wt %.\nAccordingly, it is a primary object of this invention to provide a catalytic process for upgrading a naphtha stream which contains BTX aromatics and also at least 5 wt % of C9+ aromatics, at least 25 wt % of C.sub.5 to C.sub.9 aliphatic hydrocarbons and greater than 10 wt. ppm sulfur."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of optical correlators. More specifically, the present invention discloses an optical correlator using ferroelectric liquid crystal spatial light modulators.\n2. Statement of the Problem\nOptical correlators were first suggested shortly after the advent of the laser in the early 1960's by A. Vander Lugt et al. at the University of Michigan. In early optical correlators, the input scene was introduced into the correlator by means of a photographic film transparency. The spatial filter was provided by means of a holographic film created by generating a hologram of the filter image's Fourier transform. Special care had to be taken to: (1) generate an acceptable spatial filter because of the generally large dynamic range of the Fourier transform; and (2) align the input scene's spectrum and filter encoded on the hologram. This special care translated into many hours of filter preparation and tedious mechanical alignment. However, once correlation was achieved, the classical Vander Lugt correlator did a good job of recognizing and locating patterns.\nOne major shortcoming of such simple matched spatial filtering is that the filter is extremely sensitive to differences between the object in the input scene and the object from which the filter is generated. If the difference is more than a few degrees of in-plane rotation or a few percent in scale, the focused points of light in the correlator quickly defocus and the intensity of the signal diminishes. Thus, a robust pattern recognition system for a dynamic application required a large number of filters for each input to cover the different potential orientations and scales of the target. This problem was magnified prior to the mid-1980's by the fact that the process of changing the filter meant replacing the piece of holographic film to within a few wavelengths of light. Thus, until recently, optical correlation was viewed as good physics, but was not practical for most pattern recognition applications in the field.\nA key enabling technological development in recent years is the spatial light modulator, or SLM. SLM's can be thought of as programmable transparencies or pieces of film. The use of SLM's, instead of film, in an optical correlator allows the system to rapidly change the input scene and the spatial filter without mechanically moving or replacing pads, thus accommodating the multiple-filter requirement necessary for practical pattern recognition system. SLM's have different programming speeds and resolutions depending on the type of material and the technique used to encode the scene information on the SLM. The two types of SLM's that appear best suited for use in two-dimensional pattern recognition systems are the magneto-optic SLM (MOSLM) and the family of liquid crystal SLM's.\nMOSLM devices are commercially available in pixel densities of up to 256.times.256 and have been demonstrated to operate at over 2000 Hz in short bursts, with more practical operating frame rates of 500 Hz for a 128.times.128 device and 100 Hz for a 256.times.256 device. The modulating principle of the MOSLM is Faraday rotation of the polarization vector of the incident light as the light transmits through the MOSLM. The pixels are independently and electronically addressed and are capable of binary amplitude, binary phase, or ternary phase-amplitude (combination of binary amplitude and binary phase) modulation.\nLiquid crystal technology has long been used for incoherent imaging in such applications as digital clocks, watches, and television displays. The majority of these devices use a nematic liquid crystal material that provides analog modulation, but is limited in switching speed. The most prominent nematic liquid crystal SLM for coherent imaging is the liquid crystal light valve, or LCLV. Unlike an MOSLM, an LCLV is optically addressed, rather than electrically addressed. This requires the LCLV to be programmed by another light source such as the illumination from a mini-CRT display. An LCLV uses the birefringence property of the crystalline structure and a controlled design thickness to achieve its modulation capability. The device has a maximum resolution of approximately 30 line pairs per millimeter, which equates to pixel densities on the order of 750.times.750. It has a maximum operating speed of approximately 25 to 30 Hz with the ability for modulating approximately 10 to 15 linear gray levels.\nA second type of liquid crystal SLM that has recently made significant advances in performance capabilities is the ferroelectric liquid crystal (FLC) used in the present invention. The basic performance differences between FLC SLM's and LCLV SLM's are their addressing capabilities and frame rates. FLC's can be either optically or electrically addressed, whereas LCLV's are only optically addressed. In addition, LCLV's operate at only approximately 30 Hz. Optically addressed FLC's have been operated at over 4500 Hz, and electrically addressed FLC's have been operated at approximately 10,000 Hz. The modulation principle of the FLC SLM is similar to the MOSLM in its operation (i.e., rotation of the modulator's optic axis), but the rotation is due to a reorientation of the liquid crystal molecules under an applied electric field instead of a Faraday effect. The FLC SLM is optically more efficient than the MOSLM since its axis rotation is .+-.22 degrees versus a few degrees for the MOSLM.\nA number of optical correlators and SLM's have been patented in the past, including the following:\n______________________________________ Inventor U.S. Pat. No. Issue Date ______________________________________ Yu 4,695,973 Sept. 22, 1987 Brooks 4,815,035 Mar. 21, 1989 Javidi 4,832,447 May 23, 1989 Moddel et al. 4,941,735 July 17, 1990 Juday 5,029,220 July 2, 1991 Marsh et al. 5,050,220 Sept. 17, 1991 Johnson et al. 5,073,010 Dec. 17, 1991 Capps 5,086,483 Feb. 4, 1992 Liu et al. 5,150,228 Sept. 22, 1992 Takesue et al. 5,150,229 Sept. 22, 1992 Moddel 5,177,628 Jan. 5, 1993 Moddel et al. 5,178,445 Jan. 12, 1993 Stappaerts et al. 5,221,989 June 22, 1993 ______________________________________\nYu discloses a real-time programmable optical correlator that incorporates a magneto-optic spatial light modulator (MOSLM) and a liquid crystal light valve (LCLV).\nJavidi discloses an optical correlator that employs a spatial modulator operating in a binary mode at the Fourier plane. The reference and input images are illuminated by a coherent light at the object plane of a Fourier transform lens system. An image detection device, such as a charge couple device (CCD) is placed at the Fourier plane of this Fourier transform lens system to detect the intensity of images. A thresholding network generates a binary output for each pixel of the Fourier transform interference intensity indicating whether the image intensity for that pixel is greater than the median intensity.\nJuday discloses an optical correlator for real-time tracking of the position of the retina during laser eye surgery.\nCapps discloses a hybrid optical/electronic processor in the general configuration of a Vander Lugt optical correlator with an input SLM 12, a first Fourier transform lens 16, a target SLM 14, a second Fourier transform lens 18, and an electronic processing array 20. The processing array 20 consists of a two-dimensional array of cells 40, each of which is connected to its nearest neighbors to facilitate peak detection.\nTakesue et al. disclose an optical correlator that generates pictorial patterns of a sum of two patterns of pictorial information to be compared and of a difference between the two patterns by a phase conjugate wave form. The system then transforms the pictorial patterns into first Fourier transform images, generates a pictorial pattern of a difference between an intensity distribution of the first Fourier transform images by the phase conjugate wave form, and transforms the pictorial pattern of a difference between an intensity distribution of the first Fourier transform images into second Fourier transform images. The optical correlator detects a cross-correlation peak of the two patterns of pictorial information for comparison at a high signal-to-noise ratio.\nLiu et al. disclose another example of an optical correlator using liquid crystal TV's (LCTV1 and LCTV2) to change the input and reference images in real time.\nMarsh et al. disclose an optical correlator for fingerprint identification. Two spatial light modulators 28 and 32 are employed to input the unknown fingerprint and a sequence of reference fingerprints for comparison.\nThe patents to Johnson et al., Moddel, and Moddel et al. disclose several types of optically addressable spatial light modulators incorporating ferroelectric liquid crystals.\nBrooks discloses an example of a scrolling spatial light modulator using an array of ferroelectric liquid crystal cells.\nStappaerts et al. disclose an example of a spatial light modulator using non-ferroelectric PLZT ceramic.\n3. Solution to the Problem\nNone of the prior art references uncovered in the search show an optical correlator using electrically addressable FLC spatial light modulators in the present optical configuration. The optical design is compact and facilitates easier system alignment. The system offers real-time pattern recognition at high image rates and at high resolution."}
{"text": "The power handling capability of a single capacitive element can be limited by power dissipation, voltage breakdown or, especially in the case of varactors, excessive capacitance distortion due to applied RF voltage.\nIn many resonators, it is desirable to combine multiple capacitive elements into a single Thevenin equivalent capacitor with increased power handling capability. It should be noted that capacitive elements can mean discrete capacitors, voltage variable capacitors, etched capacitors on a substrate, or combinations thereof. In high frequency resonators, it is difficult to connect several capacitors to a single discrete inductor. A popular solution is to connect the several capacitors to a distributed inductance.\nOne logical configuration for such a distributed inductor is a shorted coaxial line, as illustrated in FIG. 1. The end plate 10 short circuits the outer conductor 14 and inner conductor 12 at one end. Capacitors 16 couple the outer conductor and inner conductor at the other end.\nThe shorted coaxial resonator is advantageous in that the separation of the conductors can be as large as necessary to contain a desired number of radially connected capacitors without affecting the inductance of the resonator. The inductance of the shorted coaxial line is expressed by the following equation: EQU L=(Z.mu./2.pi.)*ln(b/a) (1)\nwhere Z is the length of the line, .mu. is the magnetic permeability of free space, b is the radius of the line's outer conductor, and a is the radius of the line's inner conductor. The inductance is a function of the ratio of the radii of the outer and inner conductors and is not dependent on the absolute diameter of the shorted coaxial line.\nAll distributed resonators exhibit resonance at a number of different frequencies. Establishing the desired resonance mode to be the dominant mode is critical in applications, such as oscillators, that are susceptible to operation at several frequencies. The desired resonance mode is a transverse magnetic (TM) wave in the axial direction of the shorted coaxial line, as illustrated in FIGS. 2a and 2b. The magnetic field lines are perpendicular (transverse) to the direction of wave propagation. The electric field lines are radially symmetric and equal in magnitude and sign in any cross-sectional plane of the resonator. Since the electric field lines are symmetric, each radial capacitor leg will receive an equal share of the resonator power.\nWhile this resonator is advantageous in certain respects, it is disadvantageous in others. Since the resonator is, by nature, a distributed circuit element, a distributed coupling technique is typically employed. Such techniques generally involve electromagnetic coupling to the resonator, such as by a coupling loop (as shown in U.S. Pat. No. 3,735,286), electrode or probe that causes an electromagnetic field to propagate into the resonant structure. Such coupling techniques are disadvantageous in certain applications that require a high degree of coupling to the Thevenin equivalent of the resonator.\nA second disadvantage in coupling to the short-circuited coaxial resonator is a difficulty in establishing a desired resonance mode. General coupling techniques can excite several different modes of resonance. One disadvantageous mode is the transverse electric (TE) mode, as illustrated in FIGS. 3a and 3b. The electric field is perpendicular (transverse) to the direction of wave propagation, and in any cross-sectional plane, the electric field does not have a radial distribution. This wave causes unequal power division of the resonator power into the capacitors.\nIn accordance with the preferred embodiment of the present invention, these drawbacks are overcome by providing a coupling port to a multiple capacitor, short circuited coax resonator. This port is defined by adding a second short circuited coax line across the end of the first. The outer conductors of the lines are interconnected. The inner conductor of the second line can be a wire, cylindrical element or reactive element such as a coil. The outer conductor of the second line can be cylindrical or a finite approximation, such as a hexagonal can, for ease of manufacturing. The inner conductors of the two lines are serially coupled and define, either along their length or at their ends, a coupling gap across which discrete circuitry can be connected. In a preferred form of the invention, the discrete circuitry is positioned in a region within the periphery of one of the inner conductors in order to provide an electromagnetic shield for the circuitry. By maintaining the radial symmetry of the resonator and coupler, the dominant resonance mode is a TM wave. The coupling port presents to the discrete circuitry a Thevenin equivalent tuned circuit composed of the sum of the capacitance of the symmetric legs in parallel with the inductance of the short circuited coaxial line.\nThe foregoing and additional features and advantages of the present invention will be more readily apparent from the following detailed description thereof, which proceeds with reference to the accompanying drawings."}
{"text": "The present invention relates to an optical head device for an optical disk apparatus for optically reading and writing information by irradiating a light spot onto a information recording surface of an optical disk (optical recording medium)\nAn optical head device for an optical disk apparatus generally includes an objective lens drive unit having an objective lens, and an optical system for receiving and transmitting a beam via the optical lens, and has such construction that the objective lens drive unit is disposed on a mounting base for an optical system block.\nThe objective lens drive unit is generally constructed of a movable unit including an objective lens, a focus coil, and a tracking coil, and a stationary unit provided with a magnetic circuit. The movable unit is supported by the stationary unit via a plurality of elastic supporting members. The elastic supporting member is enclosed at least partly by a shock-absorbing material.\nIn accordance with recent increase in speed of the optical disk apparatus, there is a need for a highly sensitive objective lens drive unit. In JP-A-2001-229555, an objective lens drive unit and an optical head device using the same in which a layout of a magnet, a focus coil and a tracking coil, which can provide a highly sensitive objective lens drive unit, is disclosed.\nReferring now to FIG. 8 and FIG. 9, an optical head device in the related art will be described. FIG. 8 is an exploded perspective view showing a construction of an objective lens drive unit of the optical head device in the related art. The objective lens drive unit shown in FIG. 8 has such construction that a coil and a magnet are disposed on both sides of an objective lens 210 respectively in tangential direction of an optical recording medium so as to be capable of driving the objective lens 210 with a high degree of sensitivity.\nAs shown in FIG. 8, the objective lens 210 is held by a lens holder 220 so as to oppose an information recording surface of an optical recording medium, not shown, for focusing an optical beam and irradiating the focused optical beam on the information recording surface of the optical recording medium to record or reproduce information.\nThe lens holder 220 shown in FIG. 8 has a coil unit 260a adhered on a left side surface S1, and a coil unit 260b adhered on a right side surface S2, which opposes the left side surface S1 in substantially parallel therewith. The lens holder 220 is formed substantially symmetrically with respect to a plane including an optical axis of the objective lens 210 and being parallel with the left and the right side surfaces S1 and S2, and is also formed substantially symmetrically with respect to a plane including the optical axis of the objective lens 210 and intersecting with the left and the right side surfaces S1 and S2. The coil units 260a and 260b are also formed symmetrically with respect to the planes described above. Therefore, the lens holder 220 and the coil units 260a and 260b adhered on both side surfaces of the lens holder 220 are integrated, and the center of gravity of a movable unit 200 including the objective lens 210 is located on the optical axis of the objective lens 210 in the lens holder 220.\nThe coil units 260a and 260b have focus coils 261a and 261b, which are wound into a rectangular shape and connected in series, respectively. The focus coils 261a and 261b are wound so as to generate the substantially same forces in the focusing direction when being energized. The coil unit 260a includes two tracking coils 262a and 262b connected in series are provided on both sides of the focus coil 261a. The coil unit 260b includes two tracking coils 263a and 263b connected in series on both sides of the focus coil 261b. The tracking coils 262a and 262b and the tracking coils 263a and 263b are connected in series. The tracking coils 262a, 262b, 263a, and 263b are wound so as to generate the substantially same forces in the tracking direction when being energized.\nThe lens holder 220 includes four wire wound coil-connecting portions 265. One of these wire wound coil-connecting portions 265 is connected to one terminal of the focus coils 261a and 261b connected in series via a leading portion (not shown) of a wire wound coil, and another wire wound coil-connecting portion 265 is connected to the other terminal of the focus coils 261a and 261b via a leading portion (not shown) of a wire wound coil.\nStill another wire wound coil-connecting portion 265 is connected to one terminal of the tracking coils 262a, 262b, 263a, and 263b connected in series via a wire wound coil leading portion (not shown), and still another wire wound coil-connecting portion 265 is connected to the other terminal of the tracking coils 262a, 262b, 263a, 263b. Each wire wound coil-connecting portion is connected to one end of each of four conductive elastic bodies 270 by welding or the like. The lens holder 220 supported by four conductive elastic bodies 270 and the coil units 260a and 260b adhered on both side surfaces of the lens holder 220 are integrated and constitute the movable unit 200 including the objective lens 200.\nThe other end of the conductive elastic body 270 is fixedly soldered to abase substrate 280. Accordingly, the movable unit 200 is cantilevered so as to be capable of moving with respect to the stationary unit including a yoke base 230, two yokes 231a and 231b, a wire base 240, two magnets 250a and 250b, and the base substrate 280.\nThe coil unit 260a is disposed in a magnetic circuit formed by a magnet 250a, which is adhered on the yoke 231a on the yoke base 230. The coil unit 260b is disposed in a magnetic circuit formed by the magnet 250b, which is adhered on the yoke 231b on the yoke base 230.\nThe coil surface of the coil unit 260a is disposed so as to face one magnetized surface 250as of the magnet 250a. The magnet 250a, which is substantially rectangular solid, is bipolarized; one in a recessed area and the other in a rectangular solid area to be fitted thereto, as shown by an image line 251a, which represents a magnetic boundary. The recessed area facing the coil unit 260a is magnetized in the S-Pole, and the rectangular solid area is magnetized in the N-Pole.\nThe coil surface of the coil unit 260b is disposed so as to face one magnetized surface 250bs of the magnet 250b. The magnet 250b, which is substantially rectangular solid, is bipolarized; one in a recessed area and the other in a rectangular solid area to be fitted thereto, as shown by an image line 251b, which represents a magnetic boundary. The recessed area facing the coil unit 260b is magnetized in the S-Pole, and the rectangular solid area is magnetized in the N-Pole. When the magnet 250a is rotated by 180° about an axis which extends in the direction of the optical axis of the objective lens 210 shown in FIG. 8, the same magnetized state as the magnet 250b is achieved.\nSubsequently, referring to FIG. 9, magnetized areas of the magnets 250a and 250b in the optical head device and the positional relation among the focus coils 261a and 261b and the tracking coils 262a, 262b, 263a and 263b will be described. FIG. 9A shows a positional relation between the magnet 250a and the coil unit 260a when viewing the lens holder 220 in the direction indicated by V in FIG. 8. In FIG. 9A, the magnet 250a is positioned on the near side with respect to the coil unit 260a. As described before, the recessed area of the magnet 250a facing the coil unit 260a is magnetized in a S-Pole and the rectangular solid area is magnetized in a N-Pole. As shown in FIG. 9A, the focus coil 261a is wound into a rectangular shape. The two tracking coil 262a and 262b are disposed on both sides of the focus coil 261a. The tracking coils 262a and 262b are also wound into rectangular shapes.\nAs a matter of convenience, four sides of the rectangular of the focus coil 261a is supposedly divided into areas of A, B, C, and D as shown in FIG. 9A, the side B of the focus coil 261a is disposed at the position facing the S-Pole of the magnet 250a, so that a magnetic flux travels in the direction from the surface of the drawing toward the near side (shown by a double-circle in the drawing). The side D opposing the side Bis disposed at the position facing the N-Pole, so that a magnetic flux travels in the direction from the near side toward the surface of the drawing (shown by a cross in the drawing). The side A and the side C of the focus coil 261a are disposed at the positions crossing over the N-Pole and the S-Pole of the magnet 250a. \nAs in the case of the focus coil 261a, when supposedly dividing the four sides of the tracking coils 262a and 262b into regions of A, B, C, and D as shown in FIG. 9A, the side A of the tracking coil 262a is disposed at the position facing the N-Pole of the magnet 250a, so that a magnetic flux travels in the direction from the near side toward the surface of the drawing (shown by a cross in the drawing). The side C opposing the side A is disposed at the position facing the S-Pole, so that a magnetic flux travels in the direction from the surface of the drawing toward the near side (shown by a double-circle in the drawing). The side B and the side D of the tracking coil 262a are disposed at the positions crossing over the N-Pole and the S-Pole of the magnet 250a. \nThe side A of the tracking coil 262b is disposed at the position facing the S-Pole of the magnet 250a, so that a magnetic flux travels in the direction from the surface of the drawing toward the near side (shown by a double-circle in the drawing). The C side opposing the A side is disposed at the position facing the N-Pole, and a magnetic flux travels in the direction from the near side toward the surface of the drawing (shown by a cross in the drawing). The side B and the side D of the tracking coil 262b are disposed at the positions crossing over the N-Pole and the S-Pole of the magnet 250a. \nThe tracking coils 262a and 262b are wound in such a manner that the direction of a current flowing along the side A of one tracking coil 262a is opposite from the direction of a current flowing along the side A of the other tracking coil 262b. \nFIG. 9B shows a positional relationship between the magnet 250b and the coil unit 260b when viewing the lens holder 220 in the direction indicated by an arrow V in FIG. 8. In FIG. 9B, the magnet 250b is positioned on the further side with respect to the coil unit 260b. As described above, the recessed area facing the coil unit 260b is magnetized in the S-Pole, and the rectangular solid area is magnetized in the N-Pole. As shown in FIG. 9B, the focus coil 261b is wound into a rectangular shape. The two tracking coil 263a and 263b are disposed on both sides of the focus coil 261b. The tracking coils 263a and 263b are also wound in a rectangular shape, respectively.\nThe side B of the focus coil 261b is disposed at the position facing the S-Pole of the magnet 250b, so that a magnetic flux travels in the direction from the near side toward the surface of the drawing (shown by a cross in the drawing). The side D opposing the side B is disposed at the position facing the N-Pole, so that a magnetic flux travels in the direction from the surface of the drawing toward the near side (shown by a double-circle in the drawing). The side A and the side C of the focus coil 261b are disposed at the positions crossing over the N-Pole and the S-Pole of the magnet 250b. \nThe side A of the tracking coil 263a is disposed at the position facing the N-Pole of the magnet 250b, so that a magnetic flux travels in the direction from the surface of the drawing toward the near side (shown by a double-circles in the drawing). The side C opposing the side A is disposed at the position facing the S-Pole, so that a magnetic flux travels from the near side toward the surface of the drawing (shown by a cross in the drawing). The side B and the side D of the tracking coil 263a are disposed at the positions crossing over the N-Pole and the S-Pole of the magnet 250b. \nThe side A of the tracking coil 263b is disposed at the position facing the S-Pole of the magnet 250b, so that a magnetic flux travels in the direction from the near side toward the surface of the drawing (shown by a cross in the drawing). The side C opposing the side A is disposed at the position facing the N-Pole, so that a magnetic flux travels in the direction from the surface of the drawing to the near side (shown by a double-circles in the drawing). The side B and the side D of the tracking coil 263b are disposed at the positions crossing over the N-Pole and the S-Pole of the magnet 250b. \nThe tracking coils 263a and 263b are wound in such a manner that the direction of a current flowing along the side A of one tracking coil 263a is opposite from the direction of a current flowing along the side A of the other tracking coil 263b. \nWhen the coil surface of the optical head device having the objective lens drive unit constructed as described above is disposed in substantially parallel with the radial direction of the optical disk, and the focus coils 261a and 261b and the tracking coils 262a, 262b, 263a, and 263b in the coil units 260a and 260b are energized, a force acting on the coil according to a Fleming's left hand rule is generated, and thus the lens holder 220 can be moved in desired directions. When the focus coils 261a and 261b are energized, a driving force for moving the side B and the side D in the focusing direction (vertical direction in FIG. 9), and when the tracking coils 262a, 262b, 263a, and 263b are energized, a driving force for moving the side A and side C in the tracking direction (lateral directions in FIG. 9) is generated.\nFor example, when a current is flown in the focus coils 261a and 261b in the directions indicated by an arrow if as shown in FIG. 9, a force Ff acting in the upward direction in the drawing is generated. Accordingly, the objective lens 210 can be moved correspondingly to the fluctuation of the surface of the optical disk. For example, the objective lens 210 can be moved by the focus coils 261a and 262a in the direction substantially vertical to the information recording surface of the optical disk to adjust a focusing position.\nAs shown in FIG. 9, when a current is flown in the tracking coils 262a, 262b, 263a, and 263b in the direction indicated by an arrow it, a force Ft acting toward the right in the drawing is generated. Accordingly, the objective lens 210 can be moved corresponding to eccentricity of the optical disk. For example, the tracking position can be adjusted by moving the objective lens 210 radially of the optical disk by the tracking coils 262a, 262b, 263a, and 263b. \nIn the objective lens drive unit of the optical head device having the construction described thus far, a driving action in a case in which part of any sides of the tracking coils 262a, 262b, 263a, and 263b are included in neutral areas n in the vicinity of the magnetic boundaries 251a and 251b, will be described. In the neutral areas n, the magnetic flux is not present, or is present but in a very low density.\nAs shown in FIG. 9A, a right portion of the side B of the tracking coil 262a and a left position of the side B of the tracking coil 262b are included in the neutral area n in the vicinity of the magnetic boundary 251a of the magnet 250. Therefore, when an attempt is made to move the movable portion 200 toward the right in FIG. 9, for example, by energizing the tracking coils 262a, 262b, 263a, and 263b, a downward force, which is generated at the right portion of the side B of the tracking coil 262a (included in the neutral area n), is smaller than the upward force, which is generated at the right portion of the side D of the tracking coil 262a. Therefore, in FIG. 9, when the tracking coil 262a is energized, an upward force Fe1, as well as force to move the lens holder 220 rightward, is generated.\nLikewise, an upward force, which is generated at the left portion of the side B of the tracking coil 262b (included in the neutral area n), is smaller than a downward force, which is generated at the left portion of the side D of the tracking coil 262b. Therefore, in FIG. 9, when the tracking coil 262b is energized, a downward force Fe2, as well as a force to move the lens holder 220 rightward, is generated. Therefore, from the coil unit 260a, a moment, which rotates the lens holder 220 clockwise when viewing in the direction indicated by an arrow V in FIG. 8, is generated.\nOn the other hand, as shown in FIG. 9B, the right portion of the side B of the tracking coil 263a and the left portion of the side B of the tracking coil 263b are included in the neutral area n in the vicinity of the magnetic boundary 251b of the magnet 250b. Therefore, for example, when an attempt is made to move the movable portion 200 toward the right in FIG. 9 by energizing the tracking coils 262a, 262b, 263a, and 263b, a downward force, which is generated at the right portion of the side B of the tracking coil 263a (included in the neutral area n), is smaller than an upward force, which is generated on the right portion of the side D of the tracking coil 263a. Therefore, in FIG. 9B, when the tracking coil 263a is energized, an upward force Fe3, as well as a force to move the lens holder 220 rightward, is generated.\nLikewise, an upward force, which is generated on the left portion of the side B of the tracking coil 263b (included in the neutral area n) is smaller than a downward force, which is generated on the left portion of the side D of the tracking coil 263b. Therefore, in FIG. 9B, when the tracking coil 263b is energized, a downward force Fe4, as well as a force to move the lens holder 220 rightward is generated. Therefore, a moment for rotating the lens holder 220 clockwise when viewing in the direction indicated by the arrow V in FIG. 8 is generated from the coil unit 260b. \nIn order to increase the sensitivity in the tracking direction, it is preferable to secure the lengths of the sides A and the sides C of the tracking coils 262a, 262b, 263a, and 263b as long as possible. However, the objective lens drive unit disclosed in the aforementioned publication, when the lengths of the sides A and the sides C of the tracking coils 262a, 262b, 263a, and 263b are increased, the sides B partly overlap the neutral areas n along the magnetic boundaries 251a and 251b. When energized for moving the lens holder 220 in the direction indicated by the arrow Ft in this state, a moment is generated by unnecessary forces Fe1 to Fe4. Therefore, there arises a problem in that a rolling about the tangential direction occurs and thus the optical axis of the objective lens 210, which is mounted to the lens holder 220 is inclined into the direction radially of the optical recording medium when an attempt is made to move the movable portion 220 only in the tracking direction by a predetermined extent.\nIn addition, there arises a problem in that when the widths of the magnets 250a and 250b in the focusing direction so as to avoid overlapping of the sides B of the tracking coils 262a, 262b, 263a, and 263b and the neutral areas n along the magnetic boundaries 251a and 251b, the thickness of the movable portion 200 must be increased, and thus reduction of size and weight of the apparatus is hindered."}
{"text": "Many natural biological molecules and their analogues, including proteins and polynucleotides, foreign substances and drugs, which are capable of influencing cell function at the sub-cellular or molecular level are preferably incorporated within the cell in order to produce their effect. For these agents the cell membrane presents a selective barrier which is impermeable to them. The complex composition of the cell membrane comprises phospholipids, glycolipids, and cholesterol, as well as intrinsic and extrinsic proteins, and its functions are influenced by cytoplasmic components which include Ca++ and other metal ions, anions, ATP, microfilaments, microtubules, enzymes, and Ca++-binding proteins, also by the extracellular glycocalyx (proteoglycans, glycose aminoglycans and glycoproteins). Interactions among structural and cytoplasmic cell components and their response to external signals make up transport processes responsible for the membrane selectivity exhibited within and among cell types.\nSuccessful delivery of agents not naturally taken up by cells into cells has also been investigated. The membrane barrier can be overcome by associating agents in complexes with lipid formulations closely resembling the lipid composition of natural cell membranes. These formulations may fuse with the cell membranes on contact, or what is more common, taken up by pynocytosis, endocytosis and/or phagocytosis. In all these processes, the associated substances are delivered in to the cells.\nLipid complexes can facilitate intracellular transfers also by overcoming charge repulsions between the cell surface, which in most cases is negatively charged. The lipids of the formulations comprise an amphipathic lipid, such as the phospholipids of cell membranes, and form various layers or aggregates such as micelles or hollow lipid vesicles (liposomes), in aqueous systems. The liposomes can be used to entrap the substance to be delivered within the liposomes; in other applications, the drug molecule of interest can be incorporated into the lipid vesicle as an intrinsic membrane component, rather than entrapped into the hollow aqueous interior, or electrostatically attached to aggregate surface. However, most phospholipids used are either zwiterionic (neutral) or negatively charged.\nAn advance in the area of intracellular delivery was the discovery that a positively charged synthetic cationic lipid, N-[1-(2,3-dioleyloxy)propyl]-N,N,N-trimethylammonium chloride (DOTMA), in the form of liposomes, or small vesicles, could interact spontaneously with DNA to form lipid-DNA complexes which are capable of adsorbing to cell membranes and being taken up by the cells either by fusion or more probably by adsorptive endocytosis, resulting in expression of the transgene [Felgner, P. L. et al. Proc. Natl. Acad. Sci., USA 84:7413-7417 (1987) and U.S. Pat. No. 4,897,355 to Eppstein, D. et al.]. Others have successfully used a DOTMA analogue, 1,2-bis(oleoyloxy)-3-(trimethylammonio)propane (DOTAP) in combination with a phospholipid to form DNA-complexing vesicles. The Lipofectin™ reagent (Bethesda Research Laboratories, Gaithersburg, Md.), an effective agent for the delivery of highly anionic polynucleotides into living tissue culture cells, comprises positively charged liposomes composed of positively charged lipid DOTMA and a neutral lipid dioleyl phosphatidyl ethanol amine (DOPE) referred to as helper lipids. These liposomes interact spontaneously with negatively charged nucleic acids to form complexes, referred to as lipoplexes. When excess of positively charged liposomes over DNA negative charges are used, the net charge on the resulting complexes is also positive. Positively charged complexes prepared in this way spontaneously attach to negatively charged cell surfaces or introduced into the cells either by adsorptive endocytosis or fuse with the plasma membrane, both processes deliver functional polynucleotide into, for example, tissue culture cells. DOTMA and DOTAP are good examples for monocationic lipids. [Illis et al. 2001, ibid.]\nMultivalent cations by themselves (including polyalkylamines, inorganic salts and complexes and dehydrating solvents) have also been shown to facilitate delivery of macromolecules into cells. In particular, multivalent cations provoke the collapse of oligo and polyanions (nucleic acids molecules, amino acid molecules and the like) to compact structural forms, and facilitate the packaging of these polyanions into viruses, their incorporation into liposomes, transfer into cells etc. [Thomas T. J. et al. Biochemistry 38:3821-3830 (1999)]. The smallest natural polycations able to compact DNA are the polyalkylamines spermidine and spermine. By attaching a hydrophobic anchor to these molecules via a linker, a new class of transfection vectors, the polycationic lipopolymers, has been developed.\nCationic lipids and cationic polymers interact electrostatically with the anionic groups of DNA (or of any other polyanionic macromolecule) forming DNA-lipid complexes (lipoplexes) or DNA-polycation complexes (polyplexes). The formation of the complex is associated with the release of counterions of the lipids or polymer, which is the thermodynamic driving force for lipoplex and polyplex spontaneous formation. The cationic lipids can be divided into four classes: (i) quaternary ammonium salt lipids (e.g. DOTMA (Lipofectin™) and DOTAP) and phosphonium/arsonium congeners; (ii) lipopolyamines; (iii) cationic lipids bearing both quaternary ammonium and polyamine moieties and (iv) amidinium, guanidinium and heterocyclic salt lipids."}
{"text": "1. Field of the Invention\nThe present invention relates to a method/system for utilizing a position sensor, such as a sensor for sensing the Y-Y position of a shift member in a vehicular transmission shifter, to determine the mode of engine fueling to be utilized to allow a shift to neutral without requiring disengagement of the vehicle master clutch. More particularly, the present invention relates to an adaptive technique utilizing a sensed rate of change in the Y-Y position of a shift member, such as a shift lever or shift finger movable in an xe2x80x9cH-typexe2x80x9d shift pattern or of an automated shift actuator, to determine if a forced torque reversal routine is necessary to allow a shift to neutral.\nThe present invention also relates to a method/system for causing engine torque to equal a value calculated to provide a zero driveline torque wherein, engine torque is changed at a rate which is a function of sensed throttle pedal position and/or engaged ratio.\n2. Description of the Prior Art\nManually operated shift levers and X-Y shifters for moving a shift member, such as a shift finger or the like, in an H-type shift pattern are well known in the prior art, as may be seen by reference to U.S. Pat. Nos. 5,481,170; 5,281,902; 4,899,607; 4,821,590; 4,784,007; 4,455,883; and 4,515,029, the disclosures of which are incorporated herein by reference.\nPosition sensors of the discrete position type for sensing discrete X-Y positions of a shift member in an H-type shift pattern are known in the prior art, as may be seen by reference to U.S. Pat. Nos. 4,817,468 and 4,718,307. Sensors providing analog or digital signals, the magnitude of which are representative of the X-X and/or Y-Y position of a shift member in an H-type shift pattern and/or shaft force, are known in the prior art, as may be seen by reference to U.S. Pat. Nos. 5,911,787; 5,743,143; 5,481,170; and; 5,950,491, the disclosures of which are incorporated herein by reference. Sensors for sensing the position of other shift members, such as pistons, ball screw members and the like are also known in the prior art.\nFully or partially automated mechanical transmission systems that, upon determining that a shift from a currently engaged ratio into neutral and than into a target ratio is desirable, will, while maintaining the vehicle master friction clutch engaged, initiate automatic fuel control to cause reduced torque across the jaw clutches to be disengaged, are also known in the prior art as may be seen by reference to U.S. Pat. Nos. 4,850,236; 5,582,558; 5,735,771; 5,775,639; and 6,015,366, the disclosures of which are incorporated herein by reference.\nThese systems include systems that attempt to fuel the engine to achieve a zero driveline torque, see U.S. Pat. No. 4,593,580, the disclosure of which is incorporated herein by reference, and systems that fuel the engine to force one or more torque reversals, see U.S. Pat. No. 4,850,236.\nSystems utilizing the engine fueling technique of forced torque reversals, especially repeated forced torque reversals, provide a highly reliable routine for allowing a shift to neutral, but may provide a somewhat less than desirable shift feel. Systems utilizing the engine fueling technique of aiming and attempting to remain at zero driveline torque provide a relatively high quality shift fuel, but did not provide desirable reliability of allowing an engaged jaw clutch to be acceptably disengaged.\nIn accordance with the present invention, an adaptive engine fueling control is provided which, based upon a sensed movement, preferably a rate of movement of a shift member, such as a shift lever in the Y-Y direction or the like, will utilize engine fueling to achieve the most appropriate (a) of aiming at a zero driveline torque, (b) or causing torque bumps to force torque reversals across the jaw clutch to be disengaged, or (c) a combination of thereof.\nIn a preferred embodiment, the above is accomplished by, upon determining that a shift into transmission neutral is required, causing the engine to be fueled to cause an output (flywheel) torque determined to cause a substantially zero driveline torque while monitoring shift member position. After a period of time, if the rate of change of the shift member position is greater than a reference value, than the desired substantially zero torque condition has been achieved (as expressed by the operator\"\"s ability to move the shift lever toward the Y-Y neutral position) and a torque bump routine to forced torque reversals is not required. If however, the shift member rate of movement is not at least equal to the reference value, than the torque bump routine may be necessary and will be initiated.\nAccordingly, it is an object of the present invention to provide an improved method/system for adaptively controlling engine fueling to allow a shift into neutral without requiring master friction clutch to remain disengaged. The adaptive engine fuel control is a function of sensed change or rate of change of a shift member position, such as, for example, the Y-Y position of a shift member movable in an H-type shift pattern.\nThis and other objects and advantages of the present invention will become apparent from a reading of the following description of the preferred embodiment taken in connection with the attached drawings."}
{"text": "The present invention is directed to a process for producing a pile in the earth with a pipe, the pipe being provided internally with a conduit, which emerges on the bottom side of the pipe, for feeding pile-forming, hardening material into the earth, the process comprising the introduction of the pipe into the earth, the removal of said pipe from the earth and, as the pipe is withdrawn, the filling of the space thus created with the pile-forming, hardening material.\nIn the prior art, numerous processes are known for the making of piles. A distinction must be drawn between reinforced and non-reinforced piles. For the production of reinforced piles, all sorts of pipe constructions are brought into the earth and, once the desired depth has been reached, the reinforcement is put in place and concrete poured or vice versa. The production of such concrete piles is a lengthy process which takes minutes to tens of minutes. Piles obtained in this way generally have a diameter greater than 20 cm. A known problem with such piles, especially with loose earth layers, is the risk of bulging as the pipe is vibrated or driven out, resulting from the difference in specific weight between the concrete and the surrounding earth. Vibration or shaking of the pipe as it is removed is necessary for the compaction of the concrete.\nOn the other hand, it is known to produce so-called foam concrete piles in the earth. A steel pipe is first placed in the soil down to the bearing layer, after which a plastics pipe may be placed therein which is subsequently filled with foam concrete, after which the steel pipe is removed. No reinforcement is used. Foam concrete of relatively low specific weight is used to combat the bulging effect.\nThe piles which are provided with reinforcements, which piles are known from the prior art, have a high bearing capacity. The structures lying thereon are of relatively heavy construction so as to bridge the large distance between the different piles of high bearing capacity. Foam concrete piles are often used to stabilize the earth. For foundation applications, these are not suitable.\nA process for producing a pile in the earth with a pipe is know from WO 93/16236. The process which is known herefrom comprises in sequence the following steps: the driving into the earth of a pipe having a closed pointed end, the lowering of a reinforcement into the driven-in pipe, the filling of the pipe driven into the earth and the withdrawal of this pipe from the earth only after the pipe has been filled, the pointed bottom end of the pipe being able to swing open so as to leave the pile-forming material behind in the earth. The pipes in question have a diameter between 150 and 300 mm.\nDE 3612437 discloses a pipe being introduced in the soil by vibration. Centrally a tube is provided. Before introduction the open end of the tube at the extrimity of the pipe is closed off by a cup shaped cover and at reachting the lowest point of introduction of the pipe followed by retraction cover 8 remains at the lower end of the hole so that the pipe becomes opened. A similar closing off and opening mechanism for concrete material to be pumped through a pipe is shown in U.S. Pat. No. 4,152,089 wherein a penetration plate 42 is pushed in advance of a pipe introduced in the soil and remains at the lower end of the bore provided.\nThe object of the present invention is to provide a process for producing a pile in the earth, which process can be realized relatively quickly. That is to say the number of maneuvers needed to arrive at the pile should be kept to an absolute minimum, thereby enabling a large number of piles to be put in place in one day, whereby hitherto used foundation methods can be substantially improved.\nThis object is achieved in the above-described process in that the conduit, as the pipe is introduced into the earth, is closed off on the bottom side of the pipe by means of a shut-off valve and is full of pile-forming material, and in that, immediately after this pipe reaches its lowest point, the shut-off valve is opened and the pile-forming material is pumped via the conduit, as the pipe is withdrawn, into the underlying space which is thereby created, and that the formed pile has a diameter less than 25 cm, such as 20 cm, and that the pumped pile-forming material contains no solid constituents of more than 16 mm in diameter.\nPreferably, the pipe is moved down to the bearing layer and not introduced into the bearing layer as is customary with driven piles and the like. The pile can also however be moved down into the bearing layer. The pile according to the invention has a relatively small diameter-and a correspondingly low bearing capacity. As an example, a value between 10 and 15 tonnes is given compared to 50-200 tonnes for conventional driven piles. In the piles according to the invention, it is not however necessary to apply a reinforcement. However, a reinforcement can very easily be provided, for example a reinforcement of fibres, for example steel fibres, mixed into the pile-forming material. Work can proceed especially quickly, because at the instant at which the pipe reaches the lowest point it is immediately withdrawn again, whereupon the opening which is left is directly filled with the pile-forming material. This pile-forming material can be any material known in the prior art, but it is essential that it compacting concrete mortar. These mortars consist of a mixture of cement, fly ash or other fine parts and water. The liquid mortar can have a strength classification which is known in the prior art, for example B15 or B25. Liquid mortar does not need to be compacted, thereby obviating the need to shake or vibrate the pipe as it is lifted, as well as preventing bulging.\nThe fact that the conduit is already full of pile-forming material as the pipe is introduced into the earth means that, as soon as the lowest point is reached, the pipe can start to be withdrawn and, at the same time, the space vacated as the pipe withdraws can start to be filled with pile-forming material. When the lowest point is reached, the pipe is therefore already full to the bottom with pile-forming, hardening material, thereby allowing an extra filling step to be omitted, so that a considerable time saving is achieved through the direct withdrawal of the pipe and the immediate filling of the space in the earth.\nIt has been shown that the pile according to the invention can be produced within 1 minute. During the first 30 seconds of this minute, the pipe is moved downwards and during the following 30 seconds the pipe is withdrawn as the cavity created is filled with the pile-forming material. Directly after the removal of the pipe, a reinforcement bar or some reinforcement might possibly be applied. The piles made with the invention are relatively short and are generally less than 15 metres in length, though piles of up to 25 meters in length are also conceivable.\nThe pipe can be introduced in any way which is known in the prior art. Preference is given to vibration and more especially to high-frequency vibration. High-frequency vibration has been shown reliably to prevent damage to surrounding structures. Moreover, it is possible to move especially quickly with the pipe through soft layers of earth.\nThe pipe can be closed off in any way during the downward motion, the shut-off valve being preferably located close to the bottom side of the pipe, so that the pipe is already filled with pile-forming material. This increases the speed of the process and at the same time prevents air having to be forced out of the pipe as the pile-forming material is introduced. This can be achieved with a valve disc which is known in the prior art and which is clamped on the front of the pipe. Once the lowest point has been reached and the pipe withdrawn, such a disc is left in the earth. It is also possible to reduce this disc to a plug which is fitted to the outlet opening of the conduit for the pile-forming material. All this is dependent upon the design of the closure for the pile-forming material. This can be a grid valve consisting of two plates sliding one over the other, each provided with openings, where in a first position these openings are situated in line and in the second position these openings cover each other to form a seal. This is an especially simple shut-off valve which can easily be operated from ground level. Other shut-off valves are also possible however. As an example, a so-called tube valve is cited. This is a shut-off valve in which part of the passage is delimited by a flexible wall surrounded by an annular chamber. Through the introduction of (air) pressure into such a chamber, the parts of the flexible wall are forced one against the other and the passage is closed. A shut-off valve of this type has the advantage over the above-described grid valve that when it is open a relatively large passageway is present. This is of benefit in connection with cleaning and installation of the above-stated plug. It is also possible to use ball valves, in which case the above-described base comprises a conical part which closes off the bottom side of the pipe to be placed in the earth. When the lowest point is reached, the base is disconnected and the pipe moves away from the cone, after which material transport can immediately take place. As indicated above, the invention allows a pile to be produced especially quickly. It is thereby possible fully to revise foundation methods. According to the invention, it is no longer proposed to use a limited number of high-load-bearing piles with a heavy structure on top, but rather it is proposed to provide a large number of piles according to the invention on a relatively small surface area. It is consequently no longer necessary to place a heavy structure on the piles, since the weight placed on a floor or the like is transmitted evenly into the earth as a result of the large number of piles. According to an advantageous embodiment of the invention, at least 5 piles per 10 m2 and according to a preferred embodiment at least one pile per m2 are present in such a structure. As a result of the fast way of making the piles, considerable savings are shown to accrue by comparison with the prior art. These savings apply to the provision of the piles according to the invention compared with piles according to the prior art. Added to this is the benefit of the lighter construction which is allowed by the use of a large number of piles.\nWith the invention, any sort of foundation is obtainable. Floors of commercial buildings are herein envisaged, though bodies of dykes, roadbeds, etc. are also envisaged.\nBecause the diameter of the pile made according to the invention is less than 25 cm, the diameter of the pipe which is placed in the earth will also be less than 25 cm. The pump for pumping the material from which the pile is made can be any pump which is known in the prior art but preferably comprises a hose pump. Apart from the reliability thereof, this has the advantage that the quantity of material which is metered per revolution is accurately known. In certain soil conditions, this is of importance. In such cases, it is essential that precisely the right quantity of material is introduced into the earth when the pipe is removed. If too little material is introduced, then a constriction arises in the pile to be made. If too much material is introduced, then an inadmissible thickening arises. By feeding the speed of withdrawal of the pipe to a computing device or the like, the pump speed is able to be controlled such that the volume which is left upon the withdrawal of the pipe is always perfectly regulated. This regulation can also, of course, work the other way round, that is to say that if the pump speed is constant the speed of withdrawal is controlled. A combination of the two is also possible. The pile-making process according to the invention is so simple that it can largely be automated. A device which drives the pipe into the earth can be accurately controlled by means of a GPS system. Since, apart from the vibration block, no other special structures are necessary, the device for introducing the pipes can be made in relatively light construction. This means that such a device can move over the building site without a large number of special measures having to be taken. It is in any event not necessary to erect on the building site heavy steel planking and the like, which planking results in hold-ups and logistical problems.\nSince the piles realised with the above-described process have a limited bearing capacity, such as 10-15 tonnes per pile, compared with the piles known in the prior art, it is possible without many special measures to test these piles for strength after the material concerned has hardened. Indeed, based on a safety factor of 2, the strength of a pile can already easily be determined if an approximately 30-tonne load is applied. A 30-tonne weight can be relatively simply applied using a mobile installation, without resorting to tie piles or large reaction weights which are used in the testing of other piles.\nThrough the performance of such tests, the safety factor can also be adjusted downwards.\nBy withdrawing the pipe from the earth at a speed greater than 0.5 m/s, such as 0.67 m/s, the likelihood of so-called bulging of the pile to be formed is reduced. If the withdrawal speed is such, then a pile of relatively constant cross-sectional size is realizable. In particular, it is herein advantageous if the speed of withdrawal of the pipe is greater than 1 m/s. The so-called bulging of the pile to be formed can be further reduced according to the invention by withdrawing the pipe from the earth without it being vibrated.\nIn order to minimize bulging on the one hand and constriction of the pile to be formed on the other hand and to obtain a pile of the most constant possible cross-sectional shape, according to the invention it is advantageous if the pump capacity and/or the speed of withdrawal of the pipe is/are controlled such that the space which is created beneath the pipe upon its withdrawal is essentially immediately filled with the pile-forming, hardening material.\nFor the purpose of facilitating the introduction of the pipe into the earth, according to the invention it is advantageous if, as the pipe is placed in the earth, water is injected at high pressure into the earth beneath the pipe such that the earth beneath the pipe is fluidized.\nAccording to the invention, the pipe can be brought into the earth by the application of a compression force to the pipe and/or by vibration of the pipe. The application of a compression force to the pipe in order to drive it into the earth is denoted indeed as xe2x80x9cpull-downxe2x80x9d. It is especially advantageous in this context if the pipe is driven into the outermost top layer of the earth by the sole application of a compression force, this for the purpose of preventing vibrations in this top layer, which is generally relatively loose."}
{"text": "The present invention relates to a golf club head, and in particular to an arrangement for a sole of the golf club head.\nConventionally, a pitching wedge or a sand wedge has a sole surface which forms only a simple uniform surface and, therefore, when hitting a golf ball, the pitching or sand wedge may bounce up on the ground or cannot strike properly deep into the sands, that is, such pitching or sand wedge may top the ball and thus it is difficult to control the ball by use of such conventional pitching or sand wedge.\nTo solve the above problems, there has been proposed a golf club head the sole surface of which has such a shape as disclosed in U.S. Patent publication No. 3,810,631, that is, in the sole surface thereof, there is formed a recess which extends along the longitudinal direction thereof from the toe side to the heel side thereof. However, the golf club head having this shape, in contrast to the previously described golf club head, tends to hit too deep into the ground or into the sands, which also makes it difficult to control the ball.\nIn view of the above circumstances, in Japanese Utility Model Publication No. 51-88356 of Showa, there is disclosed a golf club head in which there are formed recesses respectively in the two sides of the sole surface of the golf club head with the substantially central portion thereof excluded, in order to prevent the golf club head from bouncing up on the ground or entering deeply into the ground when hitting a golf ball. That is, the contact area of the sole surface of the golf club head is reduced to thereby prevent the golf club head from bouncing up on the ground and at the same time the central portion of the sole surface is left as it is to thereby prevent the golf club head from striking too deeply into the ground.\nEach of the recesses formed in the sole surface of the golf club head disclosed in the above-mentioned Japanese Utility Model Publication has a semi-circular section and, the two recesses include on their respective sole surface central portion sides wall surfaces which are formed so as to extend in mutually opposed directions toward the toe and heel portions. This prevents the golf club head from moving in the toe or heel direction when hitting a golf ball, so that a golfer using this golf club is able to strike the ball straight.\nHowever, an advanced golf player, according to conditions, must use properly a straight ball, a hook ball, and a slice ball intentionally. The above-mentioned wall surfaces receive resistance from the ground, which resistance provides an obstacle to the proper use of various ways of ball hitting."}
{"text": "This invention relates to semiconductor devices and more specifically relates to a superjunction type power MOSFET with increased avalanche energy.\nSuperjunction power MOSFETs are well known. The static and dynamic characteristics of such devices are also described in xe2x80x9cAnalysis of the Effect of Charge Imbalance on the Static and dynamic Characteristics of the Superjunction MOSFET by Proveen M. Shenoy, Anup Bhalla and Gray M. Dolay, Proceeding of the ISPSD \"\"99, pp. 99-102, June 1999.\nIn such devices, the avalanche capability, sometimes called xe2x80x9cruggednessxe2x80x9d is determined mainly by the means of preventing the turn on of the inherent parasitic bipolar transistor in a DMOS type MOSgated device. However, in the superjunction device, the concentration of the P type columns is chosen to maintain charge balance in the active area of the epitaxial silicon body material. This requirement lowers the avalanche capability of the device because the high field locates in the N type region of the epitaxial silicon layer, resulting in a higher base resistance Rb1 in the avalanche current path through the N type region and to the N+ source. Thus, in some designs, avalanche energy, that is, the amount of energy which is produced in avalanche without failure, has been as low as 50 millijoules. Attempts to increase this energy results in a reduction of the device breakdown voltage.\nIt would be desirable to increase the avalanche energy of a superjunction device without degrading the breakdown voltage.\nIn accordance with the invention, the P column dose in a superjunction device is increased to a value intentionally higher than that required for charge balance in apparent disregard of the accepted theory and design rules for superjunction devices. By doing so, the high field location moves from the N region to the P column, and, therefore, a lower Rb1 or lateral base resistance is experienced by the avalanche current through the P column to the source. Thus, the avalanche capability of the device is significantly improved (by a factor greater than 10) without degrading breakdown voltage.\nFor example, in a prior design using an N epi layer concentration of 1.26E15 and a P column dose of about 1E13, the P column dose was increased to 1.1E13 and avalanche energy was increased from 50 millijoules to 2500 millijoules. The P column dose to be used is dependent on die size and it was found that a higher dose can be used on smaller area die. Thus, a P column dose of 1.2E13 was used for a die of size 110xc3x97140 mils; while a dose of 1.1E13 was used on larger die of 257xc3x97330 mils and 315xc3x97450 mils. In all cases, a dose of 1.1E13 to 1.3E13 can be used to improve avalanche capability without adversely affecting breakdown voltage."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device with projection electrodes (bumps) which can be connected electrically and mechanically to a circuit substrate, a method of fabricating a semiconductor device for the same, and a structure for mounting the semiconductor device which electrically and mechanically connect the semiconductor device to the circuit substrate.\nFor the above-mentioned circuit substrate, a resin substrate made of a fiber glass reinforced epoxy or the like, a ceramic substrate, or a glass substrate or the like for constituting a liquid crystal display panel can be used.\n2. Description of the Related Art\nConventionally, there exists a semiconductor device with bumps which can be connected electrically and mechanically to a circuit substrate, for example a semiconductor device 71 as shown in FIG. 15.\nIn the semiconductor device 71, multiple electrode pads 74 are formed on an upper surface 72a of a semiconductor chip 72, and the upper surface 72a is covered with an insulating film 76 leaving only the upper surfaces of the electrode pads 74 not covered. Lower electrodes 79 are formed above respective electrode pads 74. Bumps 82 formed in a substantially vertical straight-walled shape with respect to the upper surface 72a of the semiconductor chip 72 are provided above respective lower electrodes 79.\nNext, a method of fabricating the semiconductor device 71 shown in FIG. 15 will be described with reference to FIG. 16 to FIG. 18 as well.\nFor fabricating the semiconductor device 71, the insulating film 76 is first formed over the entire surface of the semiconductor substrate 70 on which plural electrode pads are placed in lines in the direction perpendicular to the upper surface of the semiconductor substrate and thereby plural semiconductor chips are to be formed, as shown in FIG. 16.\nThereafter, the insulating film 76 is patterned to expose the upper surface of the electrode pads 74 with a photolithography treatment and an etching treatment.\nNext, a common electrode film 78 is formed on the entire surface of the surface 70a of the semiconductor substrate 70, including on the insulating film 76 and on the electrode pads 74 with a sputtering process.\nIncidentally, the common electrode film 78 is made by sequentially forming aluminum in a thickness of 0.8 .mu.m, chromium at 0.01 .mu.m, and copper at 0.8 .mu.m, respectively, from the side of the semiconductor substrate 70.\nThe common electrode film 78 having such a multi-layered structure serves as a connecting layer to the electrode pads 74 and a barrier layer for preventing interdiffusion, as well as an electrode used for forming bumps with a plating process.\nNext, as shown in FIG. 17, a photoresist 80 is formed over the entire surface of the common electrode film 78 in a thickness of 17 .mu.m with a spin coating process. The photoresist 80 is patterned with a photolithography treatment by using a predetermined photomask to expose the photoresist 80 in an exposure apparatus and then performing a development treatment. By the patterning, the photoresist 80 exposes the common electrode film 78 to open each area in which a bump 82 is designed to be formed later.\nNext, the bump 82, in straight-walled shape and having a thickness ranging from 10 .mu.m to 15 .mu.m, is formed in each opening of the photoresist 80 above the common electrode film 78 with a gold plating treatment which uses the common electrode film 78 as an electrode for plating.\nAfter the photoresist 80 is removed, each of the bumps 82 is used as a mask and the common electrode film 78 is etched with a wet etching process to form the lower electrodes 79 in areas aligned with the bumps 82 as shown in FIG. 18.\nFinally, the semiconductor substrate 70 is processed by cutting (dicing) at the boundary parts between the adjacent semiconductor chips of semiconductor substrate 70 into single pieces of a semiconductor chip 72 as shown in FIG. 15, thus obtaining plural semiconductor devices 71.\nIncidentally, the wet etching is performed when the common electrode film 78 is etched to form the lower electrodes 79 as described in FIG. 18 and FIG. 19 for the following reasons.\nSpecifically, since the common electrode film 78 is formed in a three-layered structure consisting of aluminum in a thickness of 0.8 .mu.m, chromium at 0.01 .mu.m, and copper at 0.8 .mu.m, respectively, from the side of the semiconductor substrate 72, the use of a dry etching process would need to use a composite etching gas as an etching gas used to obtain an etching selectivity of a layer being etched with respect to the other layers, thereby complicating the selection of the composite etching gas.\nAlso, the dry etching process is disadvantageous in industrial manufacturing since it takes a very long time for the etching. Furthermore, the dry etching process has a disadvantage in that expensive apparatuses used for the etching treatment are required.\nAccording to the wet etching process however, an etchant allowing for good etching selectivity can be selected to conveniently perform the etching treatment without requiring large-scale equipments.\nNext, an example of a conventional structure for mounting a semiconductor device to connect the semiconductor device 71 formed in accordance with the above-mentioned fabricating method to a circuit substrate will be described with reference to FIG. 19, with a liquid crystal display panel taken as an example.\nIn FIG. 19, portions other than the semiconductor device 71, glass substrates 86a, 86b of liquid crystal display panel 86, and a flexible printed circuit board (FPC) 68 are shown as a sectional view.\nIn a liquid crystal display panel 86 shown in FIG. 19, a liquid crystal 96 is filled between glass substrates 86a and 86b. A plurality of transparent electrodes 88a, 88b are provided on the opposite surface of the glass substrate 86a so as to run perpendicular with each other.\nTo mount the semiconductor device 71 above the glass substrate 86a which is the circuit substrate of the liquid crystal display panel 86, the semiconductor device 71 is disposed upside down with respect to the orientation shown in FIG. 15. The semiconductor device 71 is disposed with the bumps 82 thus located on the lower positions in alignment with the transparent electrode 88a on the glass substrate 86a.\nAt this point, an anisotropic conductive adhesive 54 is interposed between the bumps 82 and the glass substrate 86a.\nIncidentally, the anisotropic conductive adhesive 54 is formed by mixing conductive particles 52 into an insulating adhesive.\nIn this manner, while the semiconductor device 71 is set on the glass substrate 86a of the liquid crystal display panel 86, concurrently with the semiconductor device 71 being pressed on the liquid crystal panel substrate 86a, a heating is performed therefor to connect each bump 82 electrically to the respective transparent electrodes 88a on the glass substrate 86a.\nFurthermore, an FPC 68 is also disposed above the terminal electrode 88c formed on the right side in FIG. 19 on the glass substrate 86a such that an anisotropic conductive adhesive 54 having the conductive particles 52 mixed therein is interposed between the FPC 68 and the terminal electrode 88c. While the FPC 68 is pressed on the glass substrate 86a, heating is performed therefor.\nThe FPC 68 is a film patterned with a copper wiring electrode for transmitting electrical power and providing an input signal to the semiconductor device 71.\nThe above-mentioned structure holds the conductive particles 52 of the anisotropic conductive adhesive 54 between the bumps 82 and the transparent electrodes 88a and between the FPC 68 and the terminal electrode 88c, respectively, which provides an electrical connection between the bumps 82 and the transparent electrodes 88a and between the copper wiring electrode on the FPC 68 and the terminal electrode 88c as well as a mechanical connection thereof using the insulating adhesive.\nThereafter, as shown in FIG. 19, a mold resin 62 is applied onto the semiconductor device 71, the upper surface of the flexible film 68, and the periphery thereof. The mold resin 62 prevents moisture from entering each of the connections between the bumps 82 and the transparent electrodes 88a and the connection between the copper wiring electrode on the FPC 68 and the terminal electrode 88c, and also provides mechanical protection, thereby improving reliability.\nHowever, the above-mentioned conventional semiconductor device and the structure for mounting the same on the circuit substrate have a disadvantage in that a large area is occupied to mount the semiconductor device on the circuit substrate.\nThus, it is difficult to reduce the circuit substrate to which the semiconductor device is connected, in size.\nFor example, a connecting area 8 for connecting the semiconductor device 71 above the glass substrate 86a of the liquid crystal display panel 86 to the FPC 68 as shown in FIG. 19 is required to have a width of approximately 5 mm since the semiconductor device 71 has a width of 2 mm and requires a margin for connection of 1 mm and the flexible film 68 requires a margin for connection of 2 mm.\nThe portion of the glass substrate 86a at which the semiconductor device 71 and the flexible film 68 are connected is a portion which serves as a non-displaying region of the liquid crystal panel. Thus, the above-mentioned conventional structure has a disadvantage in that the area of the non-displaying region is very large with respect to the area of the displaying region.\nIn other words, a larger area is occupied to mount the semiconductor device and the FPC onto the liquid crystal display panel."}
{"text": "Some switch systems capable of operating in a remote environment are currently available. Cunningham et al. (U.S. Pat. No. 4,388,965) disclose an automatic thermal switch that controls heat flow between two thermally conductive plates. In a normally open switch configuration, the environmental temperature to which a first plate is exposed heats the plate. A rise in temperature of the first plate drives a phase-change of ammonia, Freon, or deionized water to a gas inside a power unit which is capable of driving a piston. The increased pressure caused by the phase change motivates the piston to create a thermal path between the first and second plates. The thermal path is maintained so long as the temperature of the first plate is maintained at or above the phase-change temperature of the working substance by the ambient temperature. In a normally closed switch configuration, raising the temperature of one of the thermally conductive plates causes a reaction that breaks a thermal path between the first and second plates.\nApplicants filed U.S. Non-Provisional application Ser. No. 11/467,431, on Aug. 25, 2006, entitled APPARATUS, SYSTEM, AND METHOD FOR MODIFYING A THERMAL CONNECTION, which is incorporated herein by reference in its entirety. In that application, Applicants disclosed, inter alia, a wax actuator that could extend and/or retract a plunger to motivate a thermal connector in a cryogenic atmosphere. Applicants learned after filing the application, however, that the anticipated actuator was not available and was incapable of performing the expected functions. Application Ser. No. 11/467,431 was abandoned."}
{"text": "Customer contact centers provide an important interface for customers/partners of an organization to contact the organization. The contact can be for a request for a product or service, for trouble reporting, service request, etc. The contact mechanism in a conventional call center is via a telephone, but it could be via a number of other electronic channels, including email, online chat, etc.\nThe contact center consists of a number of human agents, each assigned to a telecommunication device, such as a phone or a computer for conducting email or Internet chat sessions, that is connected to a central switch. Using these devices, the agents generally provide sales, customer service or technical support to the customers or prospective customers of a contact center, or of a contact center's clients. Conventionally, a contact center operation includes a switch system that connects callers to agents. In an inbound contact center, these switches route inbound callers to a particular agent in a contact center or, if multiple contact centers are deployed, to a particular contact center for further routing. When a call is received at a contact center (which can be physically distributed, e.g., the agents may or may not be in a single physical location), if a call is not answered immediately, the switch will typically place the caller on hold and then route the caller to the next agent that becomes available. This is sometimes referred to as placing the caller in a call queue. In conventional methods of routing inbound callers to agents, high business value calls can be subjected to a long wait while the low business value calls are often answered more promptly, possibly causing dissatisfaction on the part of the high business value caller.\nIn many call centers, the agents answering calls are organized into a plurality of groups or teams, with each group having primary responsibility of the calls in one or more call queues. Different agent groups often have responsibility for different goals or functions of the call center, such as generating customer leads, closing sales with prospects and servicing existing customers. Routing an inbound caller to an appropriate group or team of the call center to address the needs of that caller can be a burdensome, time-consuming process.\nThere is a need for a system and method for identifying high business value inbound callers at a call center during a time period in which inbound callers are awaiting connection to an agent. Additionally, there is a need to improve traditional methods of routing callers, such as “round-robin” caller routing, to improve allocation of limited call center resources to high business value inbound callers. Further, there a need to improve traditional methods of routing callers to a group or team of agents appropriate to the caller's needs from a plurality of agent groups that implement different functions or goals of the call center."}
{"text": "Welding is a well-known process that is used to join parts together for a variety of intended applications. One such welding technique is referred to as arc welding.\nArc welding is a welding technique that joins separate metal parts together using a fusion process. Intense heat is applied to metal at a joint interface between two parts, causing the metal to melt and intermix, either directly between the parts or more typically with an intermediate molten filler metal. Upon cooling and solidification, a metallurgical bond is formed between the two parts.\nIn arc welding, an electric arc produces the intense heat necessary to melt metal. The arc is formed between the workpiece and an electrode, typically in the form of a wire that is manually or automatedly guided along a weld path over the joint interface. The electrode can be in the form of a wire that not only conducts the current but also melts and supplies filler metal to the joint, also referred to as a consumable electrode.\nIn a conventional arc welding system, an arc welder cable conveys inert gas and an electrical current from suitable gas and power sources, respectively, to a welding torch. A welding wire, or electrode, can also be fed through the arc welder cable to the welding torch.\nThere are many challenges in designing an arc welder cable. In particular, when used in a robotic welding system, the welding cable must endure repeated bending and continuous operation. Robotic welding systems can include a robotic assembly that can manipulate the welding torch at a high rate of speed over a defined movement path that is repeated for each welding operation, which can make the arc welder cable more susceptible to earlier fatigue breakages. In addition, the welding cable is subjected to stretching, compressing, and bending actions performed by the robotic assembly that a human operator would not be physically capable of doing.\nIt will be appreciated that this background description has been created by the inventor to aid the reader, and is not to be taken as an indication that any of the indicated problems were themselves appreciated in the art. While the described principles can, in some aspects and embodiments, alleviate the problems inherent in other systems, it will be appreciated that the scope of the protected innovation is defined by the attached claims, and not by the ability of any disclosed feature to solve any specific problem noted herein."}
{"text": "This invention relates to a one piece portable nail care device assembly having a variety of fingernail grooming accessories stored in a covered compact casing.\nSeveral manicuring devices for fingernail grooming are known. Fingernail grooming includes trimming, shaping, filing, cuticle removing, cuticle pushing, polishing, buffing and cleaning. These devices usually includes a casing enclosing a motor for operating a detachable grooming accessory or hereinafter also referred to as bit. These accessories or bits are not housed with the motor. Only one bit can be left attached with the motor. Consequently, the other accessories or bits are left loose individually. The latter pose the problem of losing these bits which entail cost to replace. This problem may be solved by providing another casing for the accessories which now requires one to carry two casings instead of one to be able to do several desired grooming functions. The problem with this is to remember bringing both casings. An alternate approach is one proposed by U.S. Pat. No. 4,137,926 which stores the casing holding the motor inside a larger casing holding the bits. Consequently, the casing of this invention is of a size not easily carried in a ladies\"\" purse or bag especially when one already carries in it cosmetics, a telephone, a wallet, a palm, a brush and others.\nIt is therefore an object of this invention to provide a portable nail care device having the fingernail grooming accessories or bits stored in the same casing as the motor.\nIt is also an object of this invention to provide a nail care device that is capable of being manufactured at a low cost.\nIt is a further object of this invention to provide a nail care device that is simple to use and operate.\nA compact one piece nail care device assembly comprising a handle housing a motor having an output shaft and a power supply; a spindle connected to the shaft for receiving a removable grooming accessory referred to as bit having a bit head and a bit stem; a bit holder attaching on top of the handle for carrying and storing the bit; a cover for the bit holder to secure the bit inside the bit holder; a switch to trigger the device; and, means for attaching parts of the nail care device assembly together. The power to run the motor may be supplied by a regular nonrechargeable battery, a rechargeable battery or by plugging into a current supply. The spindle may be covered by a spindle cover to shield the shaft and the spindle. The bit holder may attach to the spindle cover instead of the handle when a spindle cover is used. The main feature of the claimed invention is the covered bit holder having a plurality of tubular openings attached to the handle or the spindle cover for carrying the and storing the bit."}
{"text": "The subject matter described herein relates generally to wind turbines and, more particularly, to a method and system for monitoring a wind turbine.\nGenerally, a wind turbine includes a rotor that includes a rotatable hub assembly having multiple rotor blades. The rotor blades transform wind energy into a mechanical rotational torque that drives one or more generators via the rotor. The generators are sometimes, but not always, rotationally coupled to the rotor through a gearbox. The gearbox steps up the inherently low rotational speed of the rotor for the generator to efficiently convert the rotational mechanical energy to electrical energy, which is fed into a utility grid via at least one electrical connection. Gearless direct drive wind turbines also exist. The rotor, generator, gearbox and other components are typically mounted within a housing, or nacelle, that is positioned on top of a tower.\nAt least some known wind turbines include one or more components, such as bearings, gears, and/or rotor blades that may become worn down or damaged over time. To detect such component damage, known wind turbines often include a monitoring system that measures vibrations generated by the component during an operation of the wind turbine. Such monitoring systems may be complex and/or may require significant computational resources to extract component damage information from the measured vibrations."}
{"text": "1. Field of the Invention\nThis invention relates to a flow meter. It has particular application to a flow meter that can measure volume flow of liquid in a pipe where that flow is turbulent and where the flow does not always fill the pipe. It has particular, but not exclusive, application to measuring volume flow of milk in a milking installation.\nFlow of milk the pipes of a milking parlour is notoriously difficult to measure. The flow is turbulent, rapidly varying in direction, speed and volume, and the liquid that flows in the pipe is accompanied by a large amount of air, such that the milk fills a rapidly varying fraction of the cross-section of the pipe.\n2. Summary of the Prior Art\nConventionally, the volume flow of milk in such a situation has been measured by volumetric means. Most generally, the milk is captured in a container of known volume that is repeatedly filled. The total volume flow can be determined by multiplying the number of times that the container by its volume. An example of such a system is disclosed in EP-A-0081049. Such systems can be reasonably accurate, but necessarily have a large number of moving parts. This makes them bulky, complex and expensive to manufacture. Moreover, the moving parts are difficult to clean which can lead to a build-up of contaminants. One solution to this problem is to measure electrical properties, such as conductivity, of the milk. However, such systems are of questionable accuracy, and even slight electrical leakage can cause irritation to an animal being milked. Moreover, installation of known sensors within an existing milking installation typically causes significant change to the pattern of milk flow within the installation, possibly interfering with operation of the system.\nA principal aim of this invention is to provide a flow meter that is suitable for use with milk or other turbulently flowing liquids, that requires low maintenance and does not interfere with other processes, such as milking."}
{"text": "Technical Field\nThe invention relates to a display manufacturing method and a photo alignment process.\nRelated Art\nWith the progress of technologies, the flat display apparatus has been widely applied to various fields. Especially, the liquid crystal display (LCD) apparatus, with advantages of compactness, low power consumption and non-radiation, has gradually taken the place of cathode ray tube (CRT) display apparatuses and has been applied to various kinds of electronic products, such a mobile phone, a portable multi-media apparatus, a notebook computer, a TV, a screen, etc.\nFor achieving the purpose of wide viewing angle, the manufacturer of LCD apparatuses introduces the photo alignment process to control the orientation of liquid crystal molecules, and thereby the optical performance and yield of the LCD apparatus can be enhanced. In the photo alignment process, an ultraviolet source is used to illuminate a polymer thin film (alignment film) of a substrate so that the polymer structure of the thin film's surface can generate uneven photopolymerization, photoisomerization or photocleavage, which will induce the chemical structure of the thin film's surface into a specific orientation. Thereby, the liquid crystal molecules can be arranged for achieving the photo alignment. The photo alignment light source system (also called a polarization exposure optical system) used in the photo alignment process takes the largest part of the cost among the all apparatuses. However, in the conventional art, the exposure operation is implemented by an exposure apparatus (can have one or more ultraviolet lamps) in cooperation with only one substrate stage (carrying a glass substrate) as the photo alignment light source system. Therefore, the photo alignment process squanders the tact time on the non-exposure operation, such as exchanging the substrate, aligning the substrate or adjusting the exposure angle of the substrate, so that the exposure apparatus is not used effectively and sufficiently. Hence, for increasing the production capacity, it is required to add more photo alignment light source systems and installation room in the factory building, but thus the cost of the production line and products will be increased a lot, and therefore the competitiveness of the product will be decreased.\nTherefore, it is an important subject to provide a display manufacturing method and a photo alignment process wherein the display manufacturing method can be used effectively for decreasing the tact time of the photo alignment process, and thereby the equipment cost can be decreased and the product's competitiveness can be increased."}
{"text": "Thermal inkjet printers have experienced a great deal of commercial success since their inception in the early 1980's. The fundamental principles of how thermal inkjet printers work is analogous to what happens when a pot of coffee is made. Using the electric drip coffee maker analogy, water is poured into a container (reservoir) and is channeled towards a heating element that is located at the base of the container. Once the coffee has been placed in the filter, the coffee maker is turned on and power is supplied to the heating element that is surrounded by water. As the heating element reaches a certain temperature, some of the water surrounding it changes from a liquid to a gas, thus, creating bubbles within the water. As these \"super heated\" bubbles are formed, heated water surrounding these bubbles is pushed from the reservoir into a tube and finally into the carafe. Referring now to the thermal printhead, ink is located in a reservoir that has a heating element (heater resistor) at its base. When the heater resistor is turned on for a certain amount of time (pulsed by electronic circuitry) corresponding to a certain temperature, the ink surrounding the heater resistor changes from a liquid to gas phase, thus, creating a bubble that pushes surrounding ink through an orifice and finally onto a printing medium (carafe). The aforementioned example radically simplifies inkjet technology. For a more detailed treatment of the history and fundamental principles of thermal inkjet technology, refer to the Hewlett-Packard Journal, Vol. 36, No. 5, May 1985.\nIn the coffee maker analogy, the water was poured into a container (reservoir) and channeled to a heating element located at its base. This channeling, for an inkjet cartridge, may be accomplished in a variety of different ways with the objective being to simultaneously provide the ink ejecting heater resistors with a continuous supply of ink.\nThe ink channel has traditionally been a challenging feature to fabricate both in terms of manufacturing repeatability and manufacturing cost. When manufacturing a multiplicity of printheads, variation in critical dimensions can be cataclysmic. For example, if a channel's width is too narrow, it may restrict the flow of ink to the heater resistor(s) consequently causing variations in the volume of ink ejected onto the printing medium. Likewise, if the channel width is too large, ink may be more readily supplied to some heater resistors than others thus creating variations in the rate at which ink may be ejected from the printhead nozzles (hence, the distance through which ink travels before reaching the heater resistor impacts the speed/frequency at which the printhead operates).\nIn terms of cost, traditional techniques of fabricating ink feed channels involved \"sand blasting\" holes into a substrate as disclosed in U.S. Pat. No. 5,681,764. This technique, although effective, required very specialized equipment that varied significantly from conventional IC processing thus requiring special facilities, personnel, and equipment. Consequently, there has been many efforts in the inkjet printing community to develop techniques for fabricating ink feed channels wherein the channel dimensions could be accurately controlled using standard IC manufacturing equipment and methodology. The following US patents describe such methods and techniques in an attempt to remedy the aforementioned problem.\nU.S. Pat. No. 5,308,442 illustrates a method for isotropically etching ink feed channels employing wet chemical etching. This technique incorporates standard integrated circuit (IC) photolithography and wet etch processing methodology and provides an alternative to the traditional sand blasting approach. Additionally, it provides an improvement over the sand blasting technique wherein the path through which ink flows prior to reaching the heater resistor is shortened. This technique, however, is based purely on conventional anisotropic wet chemical etching (hereafter referred to as wet etching) from the backside of the wafer/wafer substrate subsequently limiting the dimensional control of the ink feed channel. The backside of the wafer refers to the side opposite of where nozzles will be formed.\nU.S. Pat. No. 5,387,314 discloses a technique for channeling ink from a reservoir to a heater resistor by utilizing photolithography techniques with a combination of wet etching and plasma etching (a conventional gaseous etching technique hereafter referred to as dry etching). A semiconductor wafer, such as a silicon wafer, is used with a known crystallographic orientation to accommodate channels through which ink flows to the heater resistor. Such a wafer can be etched in two prominent process steps: Firstly, trenches are anisotropically etched part way into the semiconductor from the backside of the substrate. Secondly, an isotropic dry etch is used to etch from the front side (the side upon which nozzles are formed) of the substrate thus creating a channel through the substrate. The advantages of this technique as compared to that previously described in U.S. Pat. No. 5,308,442, is that the front side dry etch offers a greater degree of dimensional control. As this is well know in the semiconductor industry, isotropic wet etch processes are, in general, more variable than dry etch processes. Combining both dry and wet etch processing was a major step whereupon dimensional control of the ink feed channel was improved. However, the aforementioned process introduces an isotropic dry etch step from the front side of the wafer thus requiring the substrate above the ink feed channel to be void of active devices or signal lines.\nMany of the aforementioned challenges associated with the fabrication of ink feed channels still persist. Consequently, there remains an opportunity to develop a manufacturing process and apparatus wherein: (1) ink feed channels dimensions can be precisely controlled, (2) the distance through which ink flows before reaching the heater resistor can be minimized, (3) and the time required to form the ink feed channel is reduced."}
{"text": "1. Field of the Invention\nThe field of the invention is data processing, or, more specifically, methods, apparatus, and products for profiling an application for power consumption during execution on a compute node.\n2. Description of Related Art\nThe development of the EDVAC computer system of 1948 is often cited as the beginning of the computer era. Since that time, computer systems have evolved into extremely complicated devices. Today's computers are much more sophisticated than early systems such as the EDVAC. Computer systems typically include a combination of hardware and software components, application programs, operating systems, processors, buses, memory, input/output (‘I/O’) devices, and so on. As advances in semiconductor processing and computer architecture push the performance of the computer higher and higher, more sophisticated computer software has evolved to take advantage of the higher performance of the hardware, resulting in computer systems today that are much more powerful than just a few years ago.\nParallel computing is an area of computer technology that has experienced advances. Parallel computing is the simultaneous execution of the same task (split up and specially adapted) on multiple processors in order to obtain results faster. Parallel computing is based on the fact that the process of solving a problem usually can be divided into smaller tasks, which may be carried out simultaneously with some coordination.\nParallel computers execute applications that include both parallel algorithms and serial algorithms. A parallel algorithm can be split up to be executed a piece at a time on many different processing devices, and then put back together again at the end to get a data processing result. Some algorithms are easy to divide up into pieces. Splitting up the job of checking all of the numbers from one to a hundred thousand to see which are primes could be done, for example, by assigning a subset of the numbers to each available processor, and then putting the list of positive results back together. In this specification, the multiple processing devices that execute the algorithms of an application are referred to as ‘compute nodes.’ A parallel computer is composed of compute nodes and other processing nodes as well, including, for example, input/output (‘I/O’) nodes, and service nodes.\nParallel algorithms are valuable because it is faster to perform some kinds of large computing tasks via a parallel algorithm than it is via a serial (non-parallel) algorithm, because of the way modern processors work. It is far more difficult to construct a computer with a single fast processor than one with many slow processors with the same throughput. There are also certain theoretical limits to the potential speed of serial processors. On the other hand, every parallel algorithm has a serial part and so parallel algorithms have a saturation point. After that point adding more processors does not yield any more throughput but only increases the overhead and cost.\nParallel algorithms are designed also to optimize one more resource—the data communications requirements among the nodes of a parallel computer. There are two ways parallel processors communicate, shared memory or message passing. Shared memory processing needs additional locking for the data and imposes the overhead of additional processor and bus cycles and also serializes some portion of the algorithm.\nMessage passing processing uses high-speed data communications networks and message buffers, but this communication adds transfer overhead on the data communications networks as well as additional memory need for message buffers and latency in the data communications among nodes. Designs of parallel computers use specially designed data communications links so that the communication overhead will be small but it is the parallel algorithm that decides the volume of the traffic.\nMany data communications network architectures are used for message passing among nodes in parallel computers. Compute nodes may be organized in a network as a ‘torus’ or ‘mesh,’ for example. Also, compute nodes may be organized in a network as a tree. A torus network connects the nodes in a three-dimensional mesh with wrap around links. Every node is connected to its six neighbors through this torus network, and each node is addressed by its x,y,z coordinate in the mesh. In such a manner, a torus network lends itself to point to point operations. In a tree network, the nodes typically are organized in a binary tree arrangement: each node has a parent and two children (although some nodes may only have zero children or one child, depending on the hardware configuration). In computers that use a torus and a tree network, the two networks typically are implemented independently of one another, with separate routing circuits, separate physical links, and separate message buffers. A tree network provides high bandwidth and low latency for certain collective operations, such as, for example, an allgather, allreduce, broadcast, scatter, and so on.\nWhen processing an application, the compute nodes typically do not utilize the nodes' hardware components uniformly for each portion of the application. For example, during a portion of the application that performs a collective operation, the compute nodes typically utilize the nodes' network components that interface with the tree network but do not utilize the components that interface with the torus network. During a portion of the application that performs mathematical operations on integers, the compute nodes typically do not need to utilize the float-point units of the nodes' processors. The manner in which the nodes' hardware components are utilized to process the different portions of the application determine the overall power consumption of the nodes while executing the application. Having information on how the compute nodes consume power while executing an application may help application developers efficiently reduce the power consumption of the application, thereby conserving valuable computing resources."}
{"text": "1. Field of the Invention\nThis invention relates to Bi oxide superconductors.\n2. Description of the Related Art\nOxide superconductors have recently been considered a potential as a high transition temperature (Tc) superconductor. Among them, a so-called Bi oxide superconductor consisting essentially of Bi--(Sr,Ca)--Cu--O exhibits an extremely high Tc with comparatively good stability to external environment such as moisture. In addition, its nontoxic property incontradistinction with Tl oxide superconductors is an advantage that makes its industrial applicability still higher.\nIt is known that the Bi oxide superconductor comes in three types depending on the length of the c-axis: 24 .ANG. phase (Tc=about 8K); 30 .ANG. phase (Tc=about 80K); and 36 .ANG. phase (Tc=about 110K). These three types of Bi oxide superconductors are known to have one, two, and three planes, each plane consisting of copper and oxygen ions (Cu--O plane) per unit cell, respectively.\nTl oxide superconductors represented by a formula: Tl.sub.m Ba.sub.2 Ca.sub.n-1 Cu.sub.n O.sub.y (where m is 1 or 2; n is the number of Cu--O planes per unit cell which is 1, 2, or 3), has a crystal structure substantially identical to that of the above Bi oxide superconductors, and are reported to have a Tc of about 20K with n=1. Also, a Tc of about 90K with n=2 and a Tc of about 120K with n=3 have been reported.\nAs revealed by these examples, it has been said that the number of Cu--O planes per unit cell intimately correlates with the Tc in the oxide superconductors and that the larger the number is, the higher the Tc becomes.\nHowever, the latest research reports that the above correlation is not always strictly applied and that higher Tcs could be obtained with smaller number of Cu--O planes by changing the carrier concentration in the superconductive material. For example, it is known that TlBa.sub.2 CuO.sub.6, one of the Tl oxide superconductors, which has only one Cu--O plane per unit cell, exhibits a variation in Tc from 0 to 90K depending on the oxygen content.\nA similar phenomenon has been observed in the Bi oxide superconductors. It has been reported that Bi.sub.2 Sr.sub.2 CuO.sub.6 (24 .ANG. phase) exhibits an improvement of the Tc from 8 to about 20K by controlling the carrier concentration while substituting part of Sr for a rare earth element such as Nd. However, the maximum Tc thus obtained, being quite low compared to that of the TlBa.sub.2 CuO.sub.6 superconductor, must be improved. It has also been revealed that the Tc value can be improved on Bi oxide superconductors of 30 .ANG. phase represented by Bi.sub.2 (Sr,Ca)Cu.sub.2 O.sub.8 by controlling the carrier concentration through substitution of part of Sr and/or Ca for rare earth elements. However, since the rate of Tc improvement is small in this system, further efforts in increasing the Tc are strongly called for.\nGenerally, oxide superconductors of simpler structure can be formed more easily. To obtain a superconductor in the form of thin films by such method as a sputtering method, a simple crystal structure, particularly simple in the c-axis direction in case of Bi oxide superconductors, allows easier formation. From the viewpoint of reproducibility, the best Bi oxide superconductors is of Bi.sub.2 Sr.sub.2 CuO.sub.6 system whose lattice constant in the c-axis direction is smallest, followed by Bi oxide superconductors of Bi.sub.2 (Sr,Ca).sub.3 Cu.sub.2 O.sub.8 system. Thus, high quality, high Tc Bi oxide superconductive thin films will be prepared with satisfactory reproducibility only if the Tc of Bi.sub.2 Sr.sub.2 CuO.sub.6 and Bi.sub.2 (Sr,Ca).sub.3 Cu.sub.2 O.sub.8 are improved. From this standpoint, further improvements in Tc superconductivity of the Bi superconductors are strongly demanded.\nBy the way, to obtain a long conductor using the above described Bi oxide superconductors, attempts have been made to prepare a desired wire or tape by using a superconductive powder to fill, e.g., a metal coated tube, drawing it to a predetermined diameter and shape, and then subjecting it to a thermal treatment in an oxygen containing atmosphere. Since a material for the metal coated tube must have not only good conductivity but also high oxygen content during the preparation and thermal treatment of the Bi oxide superconductor, excellent oxygen permeability is also required. In addition, the thermal treatment at high temperatures require that the material be stabile with respect to the Bi oxide superconductor. From the above points, noble metals such as silver are often used as materials for the metal coated tube.\nHowever, the noble metals are not only expensive but also of poor mechanical properties in forming, e.g., a superconductive magnet. Copper and copper alloys which are inexpensive and have satisfactory mechanical properties are, on the other hand, susceptible to oxidation during thermal treatment in an oxygen containing atmosphere. The insufficient oxygen supply to the oxide superconductor in addition to their oxidation prevents copper and copper alloys from serving as materials suitable for obtaining the desired superconductivity.\nTo this end, there is a strong demand for development of the Bi oxide superconductors that exhibit excellent superconducting properties even at a sufficiently low oxygen content in a thermal treatment to reduce deterioration of the metal coated tube materials such as copper and copper alloys. Such a demand also applies not only to the preparation of superconductive wires and tapes but also to the preparation of superconductive elements by laminating an oxide superconductive thin film and a metal thin film one upon the other."}
{"text": "A pipe bending die is known from EP 1 591 174 B1, in which the bending head is actuated via a total of three rotational axis. A first rotational drive rotates the bending mandrel itself, a second rotational drive rotates the bending arm with the clamp die attached thereto, and a third drive finally rotates a cam disk in which a guide element runs which is connected with the clamp die and radially displaces the same during the rotation. This known construction is relatively complex and expensive as a result of the numerous mechanical elements with which the bending arm is provided. Moreover, there is a limitation in the bending freedom as a result of its large interference contour.\nEP 1 543892 B1 describes a bending apparatus in which the clamp die rests on a toothed rack which is displaceable via a gearwheel attached to the bending axis. In order to cause the forward feed motion of the clamp die towards the bending mandrel, the bending drive needs to be actuated which then causes the pivoting of the bending mandrel together with the clamp die after the forward feed motion and clamping has been completed. In this case however it is necessary to rotate the bending axis by actuating the bending drive until the forward feed motion of the clamp die occurs and the workpiece is clamped between clamp die and mandrel. Bending is only performed afterwards. The bending mandrel must be provided with a completely round configuration, with right and left bending not being easily possible.\nIn the case of EP 1 543 891 A 1, the clamp dies are disposed on a support which is moved on its part by means of a horizontally disposed worm radially towards the bending mandrel or away from the same. This known construction is very complex, large and relatively expensive and three different drives are required for actuation.\nIn the bending machine as known from EP 1 396 294 A 1, the clamp dies are disposed on guide paths on which they are moved radially by a spindle drive. Drive and spindle are integrated in the bending arm. This leads to the disadvantage of a very heavy and complex bending arm.\nWO 2007122346 A 1 describes a bending die arrangement in which a second rotary disk sits on a first rotary disk in an eccentric manner relative to its bending axis. A clamp die is rotatably attached in an eccentric manner on the second disk. The forward feed motion of the clamp die and the bending movement can be performed by a respective control of the three rotational movements and their superposition. However, a large basic rotary disk with a large interference contour and three drives are necessary again and the controlled drive of the various rotational axes is very complex.\nA rotary bending die of the kind mentioned above is known from GB 1 319 581 A, in which the bending mandrel sits on its rotational axis in a freely rotatable manner. The clamp die which is guided in a recess of a further rotary plate in the manner of a motion link in the tensioning or releasing direction is connected in an articulated manner with a swivel axis which is also disposed eccentrically to the rotational axis of the bending mandrel on a driven rotary disk which is concentric to the bending mandrel, which connection occurs via a swiveling lever which is linked to its swivel axis. By twisting the rotary disk relative to the bending mandrel, a workpiece to be bent can be clamped between clamp die and bending mandrel and thereafter be bent around the same by further rotation of the rotary disk together with the bending mandrel. The clamp die can be moved away from the bending mandrel and the workpiece can be released again by means of rotation of the rotary disk in the opposite direction after completed bending. In this case, since there is only one driven rotational axis it needs to be rotated by actuating the drive that is also responsible for the bending until the forward feed motion of the clamp die occurs and the workpiece is clamped on the bending mandrel. Bending is performed only afterwards. The bending mandrel needs to be provided with a completely round configuration. Right and left bending is not easily possible. Since the guidance of the clamp die occurs over its entire path of adjustment in a motion link recess of a rotary plate, a large rotary plate with a large interference contour in respect of a limitation of the bending freedom is necessary for this purpose."}
{"text": "Working machines such as excavators are known and include two tracks laterally spaced from each other. The tracks are attached to a chassis upon which a body is rotatably mounted about a generally vertical axis. The body includes a working implement such as a loading arm. At an end of the loading arm is a bucket. The excavator can move a load (such as earth) by rotating the body about a generally vertical axis. As the body rotates with a full bucket then the centre of gravity of the vehicle as a whole moves as the loading arm swings. Additionally, in some excavators the loading arm may swing about a king post relative to the body about a vertical axis.\nSome excavators have the track set at a fixed distance apart which is relatively large distance. Accordingly, if a load is picked up with the arm pointing forwards, i.e. parallel to the tracks, and that load is then swung sideways as the body is rotated 90 degrees, then the centre of the gravity of the machine will move towards that track which is located between the bucket and the rotational axis of the body. The spacing of the track is designed so that the machine is stable under these circumstances.\nOther excavators, typically small excavators, have the ability to vary the distance between the two tracks. Thus, for access purposes to gain access through narrow gates and the like, and for example when the excavator is loaded onto a relatively narrow transportation trailer, the tracks can be positioned in a retracted position where they are relatively close together. When the excavator is required to carry out work, such as moving loads, then the tracks are extended to a deployed (or extended) position where they are positioned further apart in the retracted position. This gives the excavator more lateral stability.\nSome excavators with variable track have a first track assembly mounted on two laterally extending beams and a second track assembly mounted on a further two laterally extending beams. Each pair of beams is slideably mounted in the chassis. The track assembly and associated beam can be extended away from the chassis to a deployed position and can be moved towards the chassis in a retracted position. In the deployed position the beams act as cantilevers and therefore must have a certain overlap with the chassis so as to ensure mechanical integrity. Because of this overlap, when the tracks are in their retracted position the left hand end of the beams of the right hand track is close to the inside edge of the left hand track and vice versa. This limits the amount of extension each beam can achieve relative to the chassis.\nFurthermore, the tracks are moved apart and together by a hydraulic actuator which is attached to the inside of each track assembly. The spacing apart of the tracks in the retracted position limits the length of the hydraulic actuator which in turn limits the amount by which the tracks can be extended.\nAn object of the present invention is to provide a working machine with improved track extension/retraction abilities."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to the fabrication of integrated circuits. More particularly, embodiments of the invention relate to a method of seasoning a chamber and depositing a low dielectric constant layer in the chamber.\n2. Description of the Related Art\nIn the fabrication of integrated circuits and semiconductor devices, materials such as oxides, e.g., carbon-doped oxides, are typically deposited on a substrate in a processing chamber, such as a deposition chamber, e.g., a chemical vapor deposition (CVD) chamber. The deposition processes typically result in deposition of some of the material on the walls and components of the deposition chamber. The material deposited on the chamber walls and components can affect the deposition rate from substrate to substrate and the uniformity of the deposition on the substrate. The material deposited on the chamber walls and components can also detach from the chamber components and create contaminating particles that can damage or destroy semiconductor devices.\nParticle contamination within the chamber is typically controlled by periodically cleaning the chamber using cleaning gases, typically fluorinated compounds, that are excited by inductively or capacitively coupled plasmas. Cleaning gases are selected based on their ability to bind the precursor gases and the deposited material formed on the chamber components in order to form volatile products which can be exhausted from the chamber, thereby cleaning the process environment of the chamber.\nOnce the chamber has been sufficiently cleaned of the process gases and the cleaning by-products have been exhausted out of the chamber, a seasoning step is performed to deposit a film onto components of the chamber forming the processing region to seal remaining contaminants therein and reduce the contamination level during processing. This step is typically carried out by depositing a seasoning film to coat the interior surfaces forming the processing region in accordance with the subsequent deposition process recipe.\nTypically, seasoning films have been deposited in deposition chambers using gas mixtures identical to the gas mixtures used in the deposition processes performed in the chamber after chamber seasoning. However, it has been found that seasoning films deposited from gas mixtures have identical components in identical ratios to the gas mixtures used in subsequent deposition processes, such as the deposition of carbon doped oxides, can crack or peel off from the interior surfaces of the deposition chamber, creating contaminating particles in the chamber.\nTherefore, there remains a need for a method of seasoning a deposition chamber in which the formation of contaminating particles from the seasoning film is minimized."}
{"text": "This invention relates to a method of producing either practically pure silica glass or a modified silica glass containing at least one metal other than silicon by using a silicon alkoxide solution, which may contain additional alkoxide(s) of the desired metal(s) to produce a modified silica glass.\nRecently, it has attracted increasing attention to produce silica glass or a \"modified silica glass,\" which contains in the glass matrix at least one metal other than silicon as exemplified by titanium-containing silica glass, without melting or fusing raw materials by using a silicon alkoxide solution or a mixed alkoxide solution containing a silicon alkoxide and at least one additional metal alkoxide. (In the following description, the term \"metal alkoxide solution\" means either a solution of silicon alkoxide or the aforementioned mixed alkoxide solution.)\nIn principle, a glass producing method of this type consists of the steps of first preparing a metal alkoxide solution which contains an organic solvent in addition to water, then allowing the metal alkoxide(s) in the solution to hydrolyze to increase the viscosity of the solution until gelation of the solution, and finally heating the gelled material to convert it into glass. In practice, there is the need of shaping the metal alkoxide solution into a desired form, i.e. intended form of the glass to be produced, such as a coating film on a solid body, independent film, fiber or lump at an intermediate stage of the aforementioned hydrolysis step while the solution exhibits a sufficiently high viscosity but has not yet gelled. As will be understood, a suitable viscosity of the metal alkoxide solution to carry out the solution-shaping operation is variable depending on the type of the shaping method, and in every case the shaping operation can successfully be carried out only in a relatively narrow range of the viscosity values of the solution. The viscosity of a metal alkoxide solution and the rate of its increase considerably vary depending on various factors such as the kind of the metal alkoxide, amount of the organic solvent and the rate of progress of the hydrolysis of the alkoxide. In general the viscosity of each metal alkoxide solution becomes higher as the amount of the organic solvent is decreased and as the hydrolysis of the alkoxide progresses. However, it is difficult to freely vary the amount of the organic solvent which is employed for the purpose of assisting the dissolution of the metal alkoxide and enhancing the miscibility of the metal alkoxide with water, and the rate of hydrolysis is determined fundamentally by the kind of the metal alkoxide. Since silicon alkoxides are generally low in the rate of hydrolysis, it takes a very long time such as several days for a metal alkoxide solution to reach a viscosity suited to the shaping operation.\nIn the described glass producing method, it is also a matter for consideration that the degree of hydrolysis of the alkoxide in the solution at the time of the shaping operation affects the manner of gelation of the shaped viscous solution and the properties of the glass obtained through a subsequent heating process. Until now, however, no technique has yet been proposed to control the viscosity of a metal alkoxide solution independently of the rate or degree of hydrolysis of the metal alkoxide in the solution.\nA more detailed description will be given with respect to the case of forming a glass coating film on a glass or ceramic substrate, or a surface of a solid body or any other form, by using a metal alkoxide solution. Coating of various substrates with glass film by this method is expected to have much application for the purpose of enhancing the chemical resistance and scratch resistance of the substrates and/or modifying the electrical characteristic or optical property of the substrate.\nIn this case the aforementioned shaping of a metal alkoxide solution is embodied in coating of a substrate with a film of the alkoxide solution which has become adequately viscous. More particularly, the coating is achieved by immersing the substrate in the metal alkoxide solution not yet gelled and withdrawing therefrom or by spraying the alkoxide solution onto a major surface of the substrate, or still alternatively by dropping the solution onto the surface of the substrate and swiftly rotating the substrate in a horizontal plane. In every case the thickness of the liquid film and, hence, the thickness of the glass coating film formed by this method depend primarily on the viscosity of the metal alkoxide solution.\nTo fully achieve the object of this coating technique, it is important to form a glass coating film which is uniform and high in the strength of its adhesion to the substrate surface and has a sufficient thickness, though this film is a very thin film usually less than one micron in thickness. For example, in the case of a colored glass coating film with the desire of a deep color, it becomes necessary to form the coating film to a considerably large thickness because it is impossible to freely increase the content of the coloring element in the metal alkoxide solution without failing to form a uniform glass coating film. As a different example, where the glass film coating technique is applied to a substrate that is rather inferior in its surface smoothness it is desirable that the glass coating film has sufficient thickness to ensure effectiveness of the coating over the entire surface area and ensure durability of the coating film. Therefore, it becomes necessary to perform the aforementioned step of forming a liquid film on the substrate when the metal alkoxide solution exhibits sufficiently high viscosity. Since the composition of the metal alkoxide solution must be determined principally in view of the chemical composition of the intended glass coating film rather than the viscosity of the solution, a sufficient increase in the viscosity of the solution should necessarily be achieved by allowing the hydrolysis of the metal alkoxide in the solution to proceed to considerable extent before coating of the substrate with this solution. In other words, it will be necessary to carry out the coating operation when the metal alkoxide solution is in a highly viscous state shortly before gelation of the solution. By using a metal alkoxide solution after the progress of the hydrolysis to such extent, certainly it is possible to coat the substrate with a relatively thick liquid film which can be expected to give a glass coating film of a sufficiently large thickness. Actually, however, heating of the liquid film formed in this way results in significant peeling and/or cracking of the resultant glass coating film and therefore fails to give a uniform coating film of glass on the surface of the substrate.\nThus, it is a matter of great inconvenience in the known methods of producing glass from a metal alkoxide solution that the viscosity of the solution can hardly be controlled independently of the progress of the hydrolysis of the metal alkoxide in the solution."}
{"text": "With the development of display technology, an oxide thin film transistor has been used widely due to advantages of high electron mobility, low fabrication temperature, and excellent homogeneity.\nAn organic light emitting diode (OLED) display is a display which emits light by itself. On basis of the driving manner, it can be divided into a passive matrix driving OLED (PMOLED) display and an active matrix driving OLED (AMOLED) display. The AMOLED display has been increasingly applied to various high performance display fields due to advantages of low cost, fast response, power saving, and a large operation temperature range."}
{"text": "1. Field of the Invention\nThe present invention relates to a computer program product, system, and method for determining linkage metadata of content of a target document to source documents.\n2. Description of the Related Art\nOften a target document needs to be compared against other source documents to determine if the target document has content matching or copied from the source documents. For instance, the source documents may comprise sensitive organizational documents which are confidential to the organization, such as a corporation or government body. However, sensitive documents often tend to propagate freely within an organization and across the organization network boundaries onto the Internet. One particular problem is when a sensitive document with proprietary, competitive or private information is leaked outside the organization without authorization.\nIn addition to the organizational and corporate setting, various services are available to determine whether a student paper has plagiarized documents in a database of papers. One such student plagiarism detection service is Turnitin (see turnitin.com).\nKnown techniques to determine whether a target document has copied content from other source documents, such as sensitive corporate documents or papers, involve comparing the content of the target document with the source document.\nOne technique for determining textual relatedness is Koala Document Fingerprinting (KDF), which has been used to determine the relatedness of computer science research documents such as technical reports, conference papers and journal articles. In a Koala search technique, the document (URL) is loaded to the KDF server as a textual representation. Then, using this text, a fingerprint of the document is generated. Finally, this fingerprint is matched against the current document fingerprint database to find related documents."}
{"text": "This invention relates to an analyzing apparatus primarily designed for the purpose of analyzing body fluids, comprising a feeding device, multiple-cell measuring device, to be called measuring unit in the following, which is constituted by a string of measuring cells linked by a channel, and a waste disposal unit for removal of the samples analyzed."}
{"text": "Organic light emitting devices (OLEDs) are comprised of several organic layers in which one of the layers is comprised of an organic material that can be made to electroluminesce by applying a voltage across the device. C. W. Tang et al., Appl. Phys. Lett. 51, 913 (1987). Certain OLEDs have been shown to have sufficient brightness, range of color and operating lifetimes for use as a practical alternative technology to LCD-based full color flat-panel displays. S. R. Forrest, P. E. Burrows and M. E. Thompson, Laser Focus World, February 1995. Since many of the thin organic films used in such devices are transparent in the visible spectral region, they allow for the realization of a completely new type of display pixel in which red (R), green (G), and blue (B) emitting OLEDs are placed in a vertically stacked geometry to provide a simple fabrication process, a small R-G-B pixel size, and a large fill factor.\nA transparent OLED (TOLED), which represents a significant step toward realizing high resolution, independently addressable stacked R-G-B pixels, was reported in U.S. Pat. No. 5,703,436, Forrest et al. This TOLED had greater than 71% transparency when turned off and emitted light from both top and bottom device surfaces with high efficiency (approaching 1% quantum efficiency) when the device was turned on. The TOLED used transparent indium tin oxide (ITO) as the hole-injecting electrode and a Mg--Ag--ITO electrode layer for electron-injection. A device was disclosed in which the ITO side of the Mg--Ag--ITO electrode layer was used as a hole-injecting contact for a second, different color-emitting OLED stacked on top of the TOLED. Each layer in the stacked OLED (SOLED) was independently addressable and emitted its own characteristic color, red or blue. This colored emission could be transmitted through the adjacently stacked transparent, independently addressable, organic layer, the transparent contacts and the glass substrate, thus allowing the device to emit any color that could be produced by varying the relative output of the red and blue color-emitting layers.\nU.S. Pat. No. 5,703,745, Forrest et al, disclosed an integrated SOLED for which both intensity and color could be independently varied and controlled with external power supplies in a color tunable display device. U.S. Pat. No. 5,703,745, thus, illustrates a principle for achieving integrated, full color pixels that provide high image resolution, which is made possible by the compact pixel size. Furthermore, relatively low cost fabrication techniques, as compared with prior art methods, may be utilized for making such devices.\nSuch devices whose structure is based upon the use of layers of organic optoelectronic materials generally rely on a common mechanism leading to optical emission. Typically, this mechanism is based upon the radiative recombination of a trapped charge. Specifically, OLEDs are comprised of at least two thin organic layers between an anode and a cathode. The material of one of these layers is specifically chosen based on the material's ability to transport holes, a \"hole transporting layer\" (HTL), and the material of the other layer is specifically selected according to its ability to transport electrons, an \"electron transporting layer\" (ETL). With such a construction, the device can be viewed as a diode with a forward bias when the potential applied to the anode is higher than the potential applied to the cathode. Under these bias conditions, the anode injects holes (positive charge carriers) into the HTL, while the cathode injects electrons into the ETL. The portion of the luminescent medium adjacent to the anode thus forms a hole injecting and transporting zone while the portion of the luminescent medium adjacent to the cathode forms an electron injecting and transporting zone. The injected holes and electrons each migrate toward the oppositely charged electrode. When an electron and hole localize on the same molecule, a Frenkel exciton is formed. These excitons are trapped in the material which has the lowest energy. Recombination of the short-lived excitons may be visualized as an electron dropping from its conduction potential to a valence band, with relaxation occurring, under certain conditions, preferentially via a photoemissive mechanism.\nThe materials that function as the ETL or HTL of an OLED may also serve as the medium in which exciton formation and electroluminescent emission occur. Such OLEDs are referred to as having a \"single heterostructure\" (SH). Alternatively, the electroluminescent material may be present in a separate emissive layer between the HTL and the ETL in what is referred to as a \"double heterostructure\" (DH).\nIn a single heterostructure OLED, either holes are injected from the HTL into the ETL where they combine with electrons to form excitons, or electrons are injected from the ETL into the HTL where they combine with holes to form excitons. Because excitons are trapped in the material having the lowest energy gap, and commonly used ETL materials generally have smaller energy gaps than commonly used HTL materials, the emissive layer of a single heterostructure device is typically the ETL. In such an OLED, the materials used for the ETL and HTL should be chosen such that holes can be injected efficiently from the HTL into the ETL. Also, the best OLEDs are believed to have good energy level alignment between the highest occupied molecular orbital (HOMO) levels of the HTL and ETL materials.\nIn a double heterostructure OLED, holes are injected from the HTL and electrons are injected from the ETL into the separate emissive layer, where the holes and electrons combine to form excitons.\nVarious compounds have been used as HTL materials or ETL materials. HTL materials mostly consist of triaryl amines in various forms which show high hole mobilities (.about.10.sup.-3 cm.sup.2 /Vs). There is somewhat more variety in the ETLs used in OLEDs. Aluminum tris(8-hydroxyquinolate) (Alq.sub.3) is the most common ETL material, and others include oxidiazol, triazol, and triazine.\nA well documented cause of OLED failure is thermally induced deformation of the organic layers (e.g. melting, crystal formation, thermal expansion, etc.). This failure mode can be seen in the studies that have been carried out with hole transporting materials, K. Naito and A. Miura, J. Phys. Chem. (1993), 97, 6240-6248; S. Tokito, H. Tanaka, A. Okada and Y. Taga. Appl. Phys. Lett. (1996), 69, (7), 878-880; Y. Shirota, T. Kobata and N. Noma, Chem. Lett. (1989), 1145-1148; T. Noda, I. Imae, N. Noma and Y. Shirota, Adv. Mater. (1997), 9, No. 3; E. Han, L. Do, M. Fujihira, H. Inada and Y. Shirota, J. Appl. Phys. (1996), 80, (6) 3297-701; T. Noda, H. Ogawa, N. Noma and Y. Shirota, Appl. Phys. Lett. (1997), 70, (6), 699-701; S. Van Slyke, C. Chen and C. Tang, Appl. Phys. Lett. (1996), 69, 15, 2160-2162; and U.S. Pat. No. 5,061,569.\nOrganic materials that are present as a glass, as opposed to a crystalline or polycrystalline form, are desirable for use in the organic layers of an OLED, since glasses are capable of providing higher transparency as well as producing superior overall charge carrier characteristics as compared with the polycrystalline materials that are typically produced when thin films of the crystalline form of the materials are prepared. However, thermally induced deformation of the organic layers may lead to catastrophic and irreversible failure of the OLED if a glassy organic layer is heated above its T.sub.g. In addition, thermally induced deformation of a glassy organic layer may occur at temperatures lower than T.sub.g, and the rate of such deformation may be dependent on the difference between the temperature at which the deformation occurs and T.sub.g. Consequently, the lifetime of an OLED may be dependent on the T.sub.g of the organic layers even if the device is not heated above T.sub.g. As a result, there is a need for organic materials having a high T.sub.g that can be used in the organic layers of an OLED.\nThe most common hole transporting material used in the HTL of OLEDs is a biphenyl bridged diamine, N,N'-diphenyl-N,N'-bis(3-methylphenyl)-1,1-biphenyl-4,4'-diamine (TPD) having the chemical structure: ##STR1##\nThis material has a good hole mobility and efficiently transfers holes to aluminum tris (8-hydroxyquinoline) in a simple single heterostructure OLED. However, TPD has a melting point of 167.degree. C. and a glass transition temperature of 65.degree. C. If a device prepared with TPD is heated above 65.degree. C., the glass transition temperature, catastrophic and irreversible failure results. In order to increase the glass transition temperature of the HTL, several groups have explored different modifications to the basic structure of TPD, Naito et al.; Tokito et al.; Shirota et al.; Noda et al. (Adv. Mater.); Han et al.; Noda et al. (Appl. Phys. Lett.); Van Slyke et al.; and U.S. Pat. No. 5,061,569. While these studies have led to materials with T.sub.g values as high as 150.degree. C., they have not led to an understanding of why certain structural modifications increase T.sub.g, while other modifications may not affect T.sub.g at all or may even lower T.sub.g. Still other modifications may produce a material not having a glass transition temperature at all or a material not having the combination of properties that is suitable for use in an HTL. For example, replacing the amine groups of TPD with carbazole groups to produce 4,4'-di(N-carbazolo)diphenyl (CBP), having the chemical structure: ##STR2## increases the melting point to 285.degree. C. However, the material shows no glass transition. Further changes in the basic structure of TPD can increase the T.sub.g value even higher, but the materials often have poorer hole transporting properties than TPD, i.e. OLEDs made with these high temperature materials give poor device properties in OLEDs compared to TPD.\nU.S. Pat. No. 5,061,569 discloses hole transporting materials comprised of at least two tertiary amine moieties and further including an aromatic moiety containing at least two fused aromatic rings attached to the tertiary amine nitrogen atoms. Out of the large number of compounds encompassed by the broadly disclosed class of compounds recited, U.S. Pat. No. 5,061,569 fails to disclose how to select those compounds which have a high glass transition temperature. For example, the naphthyl derivatives do make stable glasses. One such molecule is 4,4'-bis[N-(1-naphthyl)-N-phenyl-amino]biphenyl (.alpha.-NPD), having the chemical structure: ##STR3##\nThe present inventors' measurements show that .alpha.-NPD has a T.sub.g of 100-105.degree. C., which is substantially higher than the T.sub.g of 65.degree. C. of TPD. This material has excellent hole conduction properties, and the T.sub.g of 100-105.degree. C. is higher than the T.sub.g of TPD of about 65.degree. C. OLEDs prepared with NPD have electrical properties very similar to those prepared with TPD. However, 4,4'-bis[N-(2-naphthyl)-N-phenyl-amino]biphenyl (.beta.-NPD), having the structure: ##STR4## has been generally understood to have a T.sub.g which is substantially lower than .alpha.-derivative. Apparently because of this purportedly low and anomalous difference between T.sub.g of the .alpha.- .beta.-derivatives, there had been no known reports of using the .beta.-derivative as the hole transporting material of an OLED.\nIt would be desirable if OLED's could be fabricated from glassy charge carrier materials having improved temperature stability, while still providing luminescent characteristics comparable to prior art compounds. As used herein, the term \"charge carrier layer\" may refer to the hole transporting layer, the electron transporting layer or the separate emissive layer of an OLED having a double heterostructure. In addition, it would be useful to have a method for selecting and preparing such glassy charge carrier materials having improved temperature stability, as characterized, in particular, by glassy charge carrier materials having a high glass transition temperature.\nIn addition, there is a general inverse correlation between the T.sub.g and the hole transporting properties of a material, i.e., materials having a high T.sub.g generally have poor hole transporting properties. Using an HTL with good hole transporting properties leads to an OLED having desirable properties such as higher quantum efficiency, lower resistance across the OLED, higher power quantum efficiency, and higher luminance. There is therefore a need for a HTL having a high hole mobility and a high glass transition temperature."}
{"text": "Implanted medical devices or biomaterials can present microenvironments that are conducive to bacterial colonization or infection. Such infections can require repetitive hospital visits or can even require removal or replacement of the implanted device or material. For example, approximately 1 million inguinal and 20,000 ventral herniorrhaphies are performed yearly in the United States. It is estimated that approximately 3 to 8 percent of patients who undergo such surgery develop post-operative infections. These types of infections can cause complications that require longer post-surgery care and escalate health care costs. It has been estimated that surgical site infections have increased hospital stays by approximately 7 days with costs increasing from $3,000 to $30,000. Overall surgical implant infections costs have been credited with increasing direct medical costs by more than $3 billion annually in the US. In addition to the financial burden caused by repetitive implant surgery and longer inpatient hospital stays, infections can be fatal or near fatal to a significant portion of patients.\nAt least some degree of bacterial entry into a surgical site is nearly unavoidable at the time of surgery. Bacteria can be introduced in surgeries that include the implantation of medical devices and materials as well as surgeries that do not. However, the presence of an implant can complicate the clearance of bacteria from the surgical site. Bacteria introduced at the time of implantation can 30 contact the implant, bind to and colonize the implant surface. If bacterial colonization occurs on the surface of the implant, traditional means of treatment including the use of systemic antibiotics, are largely unsuccessful in treating the infection. Because of the difficulty of clearing bacteria that have populated the implant surface, it is common for surgeons to completely remove the implant, place patients on long term systemic antibiotics to clear the infection, and re-implant the medical device.\nCurrently available procedures for treating infections typically involve administration of antibiotics. Others such as Gristina et al. (U.S. Pat. No. 5,292,513) disclose methods for nonspecific cellular immune stimulation. Ziegler et al. (U.S. Pat. No. 7,906,132) disclose anti-infectious, biocompatible titanium coating for implants, and methods for the production thereof. O'Hagan et al. U.S. Pat. No. 7,597,908 discloses use of micro-particles with adsorbed antigen to stimulate immune responses towards the specific antigen at a systemic level."}
{"text": "Numerous apparatuses have been proposed for preventing overfilling of fuel tanks. The most common used method is an automatic cut-off within a nozzle. Typically this automatic cut-off uses a vapor path from the nozzle outlet back to a venturi around the fuel flow path within the nozzle. A sufficiently high pressure must be maintained at a point within this path to indicate that vapor is being drawn into the vapor path rather than liquids. When liquids enter the vapor path, the pressure drop in the path increases, and the pressure at the sensor point will decrease. When this pressure decreased below a threshold pressure, the fuel flow is cut-off, usually by a mechanical trip. When a vapor recovery system that seals the fuel inlet is incorporated in a fuel dispenser, this automatic fuel cut-off will not function properly because pressure at the pressure sensor is subject to variations due to variations in the vapor recovery system. Such variations result in the shutoff not having sufficient consistency.\nAn electo-mechanical fuel cut-off switch is disclosed in U.S. Pat. No. 5,131,441. This switch includes an electromagnetic clutch that enables the trigger of a nozzle to close. When fluids are detected by a fluid actuated switch located in the nozzle spout, electrical energy to the electromagnetic clutch is interrupted, and the fuel valve is closed. This mechanism is said to be quick-acting, and therefore minimizes splash-back losses to the environment. An optical liquid sensor is suggested as the fluid actuated switch. The fluid actuated switch is located within the nozzle in the fuel dispenser of '441. The location of the switch within the nozzle relies on the fluid level raising within the nozzle. Because vapor is trapped within the closed volume of the nozzle, liquid will not necessarily back-up in the nozzle, but could raise outside the nozzle and be drawn into the vapor recovery system. Liquids could therefore be exiting the fuel tank into the vapor recovery system before the fuel flow is cut off by the mechanism of patent '441.\nIt is therefore an object of the present invention to provide a method and apparatus for cutting off fuel flow in a fuel dispensing nozzle wherein a vapor recovery system having a seal around a vehicle's fuel inlet can be utilized without effecting the fuel cut-off."}
{"text": "(1) FIELD OF THE INVENTION\nThe present invention relates to a driving method and a driving device for selectively driving some of a plurality of light shutter elements together constituting a light shutter array, each of the elements being formed of such material as PLZT having an electro-optical effect.\n(2) DESCRIPTION OF THE PRIOR ART\nA known light shutter array of the above-noted type using such material as PLZT having an electro-optical effect, comprises a polarizer and an analyzer which are respectively constituted by a polarizing plate and are disposed at a light incident side and at a light emitting side with polarizing angles thereof crossing at a right angle to each other. For a light transmission therethrough, a driving voltage is applied to electrodes disposed at opposed sides of the shutter array, such that the light incident to the shutter through the polarizer undergoes 90 degrees polarization in the shutter so as to be transmitted through the analyzer.\nFurther, in order to increase a transmission light intensity difference between the ON state and OFF state of the shutter, a half-wavelength voltage for polarizing the light by 90 degrees is employed as the shutter driving voltage applied between the electrodes in accordance with the 90 degree polarizing angle difference between the polarizer and the analyzer.\nThen, the inventors of the present invention thought of using such a light shutter array including a plurality of serially aligned light shutter elements in an electrophotographic printer and have conducted intensive research on this application.\nIn using such a light shutter array in an electrophotographic printer, as shown in FIG. 1, a polarizer 101 and an analyzer 102 are disposed respectively at the light-incident side and at the light-emitting side across the light shutter array 10. A light from a light source lamp 103 and its reflecting mirror 104 is guided through the polarizer 101 to the shutter array 10, in which some light shutter elements are driven depending on the image information so as to permit selective transmission of this light through the analyzer 102. This light through the analyzer 102 is converged by a converging rod lens array 105 to be irradiated onto a periphery of a photosensitive member 107 uniformly charged by a corona charger 106, whereby an electrostatic latent image is formed on the periphery. Then, this latent image is developed through a developing device 108 into a toner image, which image is transferred through a transfer device 109 onto a copy sheet while removing excess toner remaining on the peripheral surface of the photosensitive member 107 by an eraser 110 and a cleaner 111.\nIn the above-described construction of the light shutter array employed in the electrophotographic printer, the respective shutter elements need to be selectively driven by pulses in order to form the electrostatic latent image on the photosensitive member being rotated at a predetermined peripheral speed.\nThus, the inventors conducted research on ambient temperature variations in a transmission light intensity Ip under the ON state of the shutter and a leakage light intensity Id under the OFF state of the same. As shown in FIG. 2, the research revealed that the transmission light intensity Ip is hardly affected by any ambient temperature variations whereas the leakage light intensity Id increases with a decrease in the temperature and further that this leakage light intensity Id under the OFF state of the shutter significantly increases as shown in FIG. 3 if the light shutter is driven at a temperature lower than the room temperature.\nMoreover, the research revealed, as indicated by a broken line in FIG. 4, that the leakage light intensity also increases due to generation of polarization in the respective shutter elements if the shutter array is continuously driven by constantly applying a unidirectional pulsate electric field thereto even if the ambient temperature is maintained constant."}
{"text": "This invention relates to the construction of a blade wherein blade cooling is desired. Many constructions of cooled rotor blades with various constructions of leading edges appear in the prior art and some patents showing different leading edge constructions are U.S. Pat. Nos. 3,606,572; 3,240,468; 3,647,316; and 3,656,863."}
{"text": "Synchronization is critical for telecommunication system performance. Frequency and time (time-of-day or wall-clock) synchronization is crucial for mobile wireless networks because the radios used in these networks operate in very strict bands that need separation to avoid channel interference which reduces the call quality and network capacity. Poor synchronization has also negative impact for the handover between base stations.\nMobile handsets generally derive the accurate frequency that they transmit and receive from the base stations. If the transmission frequencies are not very closely matched between adjacent cell sites, then “clicks” can occur when the call is being handed over (that is, switches) between base stations. In the worst case, the call would drop because the mobile handset would not be able to immediately lock onto and acquire the new signal.\nFailure to meet the timing requirements of the relevant standards would cause performance degradation for the radio access channels. In particular, this failure could compromise cell handover (especially for travelling mobile stations) and producing excess of dropped calls.\nWith increasing interest in packet networks as a common mode of communication, packet-based synchronization solutions are in high demand as alternative to PDH (plesiochronous digital hierarchy)/SDH (synchronous digital hierarchy) and GPS based solutions. Equipment vendors and telecom service providers are looking for new packet synchronization solutions with very high accuracies beyond those attainable using the traditional packet methods like Network Time Protocol (NTP) (see Mills, D., “Network Time Protocol (Version 3) Specification, Implementation and Analysis”, IETF RFC 1305, March 1992).\nIt is desirable that such new solution are also be designed for applications (e.g., base stations) that cannot bear the cost of a GPS receiver at each node, or for which GPS signals are inaccessible (for example due to location).\nThe IEEE Standard 1588 Precision Time Protocol (PTP) (IEEE Standard for a Precision Clock Synchronization Protocol for Networked Measurement and Control Systems, IEEE 1588-2008), is the latest addition in packet timing technology. Originally designed to provide precise timing for critical industrial automation applications it is now providing the highest level of accurate frequency and time to wireless backhaul networks. The backhaul portion of the network comprises the intermediate links between the core network, or backbone, of the network and the base stations. Currently standardized in 2008 as Version 2 (IEEE-1588v2, superseding IEEE 1588-2002), PTP is designed to overcome the Network Time Protocol (NTP) latency and jitter issues, providing accuracy in the nanosecond range.\nBase stations have traditionally met synchronization requirements by locking their internal oscillators to a recovered clock from the T1/E1 TDM (time division multiplexing) backhaul connection. While Ethernet has proven to be a ubiquitous and inexpensive medium for connectivity, it has not been well-suited for applications requiring precise synchronization. When the backhaul transitions from TDM to Ethernet, the base station becomes isolated from its traditional network sync feed. New base station designs are incorporating IEEE 1588 PTP slave clocks to meet the 50 ppb (parts per billion) accuracy requirement. These PTP slave clocks in the base stations rely on access to a PTP grandmaster clock deployed in a mobile switching center (MSC). Sync and Follow_Up packets flow from the grandmaster clock to the slave clocks in the base stations.\nIEEE 1588v2 PTP is fully compatible with all Ethernet and IP networks. Additionally, the protocol is designed to enable a properly designed network to deliver frequency and phase or time with precision rivalling a GPS receiver. An IEEE 1588v2 PTP Technique for Frequency Synchronization implementation can supply FDD (frequency division duplexing) and TDD (time division duplexing) radio systems and CES-based (circuit emulation services) transport systems with the synchronization signals they require as illustrated in FIG. 1. This greatly reduces the costs of clocking all wireless base station equipment using other means.\nA primary reference is a source of time and or frequency that is traceable to international standards. A recognized standard time source is a source external to PTP that provides time and/or frequency as appropriate that is traceable to the international standards laboratories maintaining clocks that form the basis for the International Atomic Time (TAI) and Universal Coordinated Time (UTC) timescales. Examples of these are Global Positioning System (GPS), NTP, and National Institute of Standards and Technology (NIST) timeservers.\nAlthough IEEE 1588v2 PTP systems add a small amount of additional traffic to the network load, they have several advantages. First, they work in the data path, and also benefit from the redundancy and resiliency mechanisms of the network, resulting in “always on” operation. Next, multiple transmission paths reduce redundant clock system costs. They also use a single synchronization session for all base station traffic. IEEE 1588v2 PTP supports any generic packet-based transport (such as IP, MPLS). The protocol also features configurable synchronization packet rates for network conditions to maintain accuracy.\nThe transmission of the clock information over a packet network eliminates the need for alternative mechanisms, such as GPS or prohibitively expensive oscillators placed at the receiving nodes. This provides significant cost savings in network equipment as well as in ongoing installation and maintenance. This synchronization solution transmits dedicated timing packets, which flow along the same paths with the data packets, reducing the cost of synchronization and simplifying implementation.\nTo ensure that packet technologies (Ethernet, IP, MPLS) have the necessary attributes to be truly carrier grade, operators and vendors are introducing several key technologies for the transport of timing and synchronization over packet networks. Of these IEEE 1588v2 PTP is perhaps the most important because it provides both frequency and time distribution.\nThe grandmaster is the root timing reference in a domain and transmits synchronization information to the clocks residing in its domain. In IEEE 1588v2 PTP messages are categorized into event and general messages. All IEEE 1588 PTP messages have a common header as shown in FIG. 2.\nEvent messages are timed messages in that an accurate timestamp is generated at both transmission and receipt of each message. Event messages have to be accurately timestamped since the accuracy in transmission and receipt timestamps directly affects clock distribution accuracy. A timestamp event is generated at the time of transmission and reception of any event message. General messages are not required to be timestamped. The set of event messages consists of Sync, Delay_Req (both of which have the format shown in FIG. 4), Pdelay_Req, and Pdelay_Resp. The set of general messages consists of Announce (which has the format shown in FIG. 3), Follow_Up (which has the format shown in FIG. 5), Delay_Resp (which has the format shown in FIG. 6), Pdelay_Resp_Follow_Up, Management, and Signalling.\nThe Sync, Delay_Req, Follow_Up, and Delay_Resp messages are used to generate and communicate the timing information needed to synchronize ordinary and boundary clocks using the delay request-response mechanism. A Sync message is transmitted by a master to its slaves and either contains the exact time of its transmission or is followed by a Follow_Up message containing this time. In a two-step ordinary or boundary clock, the Follow_Up message communicates the value of the departure timestamp for a particular Sync message.\nA Delay_Req message is a request for the receiving node to return the time at which the Delay_Req message was received, using a Delay_Resp message.\nThe format of the PTP message timestamp fields is shown in FIG. 7. A timestamp is the current time of an event that is recorded by a device. IEEE 1588 PTP allows for two different types of timestamping methods, either one step or two-step. One-step clocks update time information within event messages (Sync and Delay-Req) on-the-fly, while two-step clocks convey the precise timestamps of packets in general messages (Follow_Up and Delay-Resp).\nThe basic pattern of synchronization message exchange is illustrated in FIG. 8. The message exchange pattern is as follows. The master 101 sends a Sync message (M101) to the slave 102 and notes the time T1 at which it was sent. The slave 102 receives the Sync message and notes the time of reception T2. The master 101 conveys to the slave 102 the timestamp T1 by one of two ways: 1) Embedding the timestamp T1 in the Sync message. This requires some sort of hardware processing (i.e., hardware timestamping) for highest accuracy and precision. 2) Embedding the timestamp T1 in a Follow_Up message (M102). Next, the slave 102 sends a Delay_Req message (M103) to the master 101 and notes the time T3 at which it was sent. The master 101 receives the Delay_Req message (M103) and notes the time of reception T4. The master 101 conveys to the slave 102 the timestamp T4 by embedding it in a Delay_Resp message (M104).\nAt the end of this PTP message exchange, the slave possesses all four timestamps {T1, T2, T3, T4}. These timestamps may be used to compute the offset of the slave's clock with respect to the master and the mean propagation time of messages between the two clocks. The computation of offset and propagation time assumes that the master-to-slave and slave-to-master propagation times are equal.\nThe IEEE 1588 PTP based frequency recovery technique described in Section 12.1 of IEEE 1588-2008 Standard requires estimation of the mean path delay between server (master) and client (slave) which may include accounting for path asymmetry. In order to accurately synchronize to their master, slave clocks must individually determine the network transit time of the PTP messages. The network transit time is determined indirectly by measuring round-trip time from each slave to its master.\nLike all message-based time transfer protocols, PTP time accuracy is degraded by asymmetry in the paths taken by event messages. Any asymmetry in the forward and reverse path propagation times and introduces an error into the computed value of the link delay. Asymmetry can be introduced in the physical layer, e.g., via transmission media asymmetry, by bridges and routers, and in large systems by the forward and reverse paths traversed by event messages taking different routes through the network. Systems should be configured and components selected to minimize these effects guided by the required timing accuracy. In single subnet systems with distances of a few meters, asymmetry is not usually a concern for time accuracies above a few 10 s of ns. Asymmetry is not detectable by PTP; however, if known, PTP corrects for asymmetry. If two-step clocks are used, then the network has to be designed such that the general message takes the same path as the event message through a transparent clock. Failure to do this will result in a condition where the transparent clock does not calculate path delay properly. This condition is undetectable and may introduce additional jitter and wander, but it will not break the protocol.\nOther IEEE 1588 based techniques (such as those discussed in T. Neagoe and M. Hamdi, “A Hardware IEEE-1588 Implementation with Processor Frequency Control,” Arrow Electronics White Paper; T. Neagoe M. Hamdi and V. Cristea, “Frequency Compensated, Hardware IEEE-1588 Implementation,” IEEE International Symposium on Industrial Electronics, 9-13 Jul. 2006, pp. 240-245 and S. Balasubramanian, K. R. Harris, A. Moldovansky, “A Frequency Compensate Clock for Precision Synchronization Using IEEE1588 and its Applications to Ethernet” IEEE-1588 Workshop, September 2003 assume that the PTP GrandMaster Clock (GMC) sends Sync messages at fixed intervals, an assumption which might not necessary hold true in practice. The above receiver clock recovery mechanisms are designed based on this assumption.\nAn object of the present invention is to provide a method and system for frequency synchronization that allows one or more receivers (slaves) to frequency synchronize to a transmitter (master). Other applications of frequency synchronization are in process and manufacturing industries like paper mills, printing presses, automation and robotic systems, test and measurement instruments and systems, etc.\nAnother object of the present invention is to provide a synchronization technique based on, for example, the IEEE 1588 Precision Time Protocol (PTP) that allows frequency to be distributed over a packet network from a PTP server (master) to synchronization clients (slaves)."}
{"text": "The differentiated cell state is often considered stable and resistant to changes in lineage identity. However, differentiated somatic cell types from humans and other organisms have been reprogrammed to the pluripotent state (“pluripotent reprogramming”) by forced expression of a set of transcription factors (Takahashi, K. et al. Induction of pluripotent stem cells from adult human fibroblasts by defined factors. Cell 131, 861-872 (2007)), somatic cell nuclear transfer (Campbell et al., Sheep cloned by nuclear transfer from a cultured cell line. Nature 380: 64-66 (1996); Gurdon et al., Sexually mature individuals of Xenopus laevis from the transplantation of single somatic nuclei, Nature 182, 64-65 (1958)) or cell fusion (Cowan et al., Nuclear reprogramming of somatic cells after fusion with human embryonic stem cells, Science (New York, N.7309, 1369-1373 (2005); Tada et al., Nuclear reprogramming of somatic cells by in vitro hybridization with ES cells, Curr Biol 11, 1553-1558 (2001)). Additionally, a few studies have demonstrated that through ectopic expression of selected genes or by cell fusion, an adult cell type can be directly converted to another adult cell type (Cobaleda et al., Conversion of mature B cells into T cells by dedifferentiation to uncommitted progenitors, Nature 449, 473-477 (2007); Davis et al., Expression of a single transjected cDNA converts fibroblasts to myoblasts, Cell 51, 987-1000 (1987); Feng, et al. PU.l and C/EBPalpha/beta convert fibroblasts into macrophage-like cells, Proc. Nat. Acad. Sci. USA 105, 6057-6062 (2008); Ieda et al. Direct reprogramming of fibroblasts into functional cardiamyocytes by defined factors, Cell 142, 375-386 (2010); Zhou et al., In vivo reprogramming of adult pancreatic exocrine cells to beta-cells, Nature 455, 627-632 (2008); and Zhou, Q. & Melton, D. A. Extreme makeover: converting one cell into another, Cell Stem Cell 3, 382-388 (2008)). This process is termed trans-differentiation or lineage reprogramming.\nHowever, major challenges remain due to the low efficiency and slow reprogramming process. A more significant challenge is how to accomplish cell reprogramming without the need for genetic changes in the reprogrammed cells, because such genetic changes give rise to concerns about introduced mutations at the insertion site of expression cassettes encoding pluripotency factors."}
{"text": "Most catheters are a cannula or tube like device which is inserted into a portion of a person's body in order to transport fluids or gases in or out of that particular portion of the body. In passing through any particular portion of the body in order to reach its destination, the catheter will come into contact with various tissues in the body. For example, a catheter used to drain one'bladder (such as a \"Foley\" catheter) must pass through the urethral tract in order to reach the bladder. A nasogastric catheter must pass through the nasal passageway and the esophagus in order to reach the stomach. Some catheters, such as these, are inserted through existing passageways in order to reach their destinations, while others are inserted through surgically created passageways.\nIn virtually every catheterization, there is a significant potential for microbial growth along the exterior surface of the catheter which can lead to serious infections such as urinary tract infections, bladder infections and the like. Such an infection can be encouraged when adjacent tissues are inflamed due to irritation from rubbing a chafing against this catheter. This is because inflamed or irritated tissues may be less apt to respond effectively to suppress local bacterial infection. In such a situation the infection can spread and intensify, placing the patient at further risk. Such infections can lead to sepsis of the bladder particularly in elderly patients who are incontinent and have a chronic need for catheterization with an indwelling self-retaining catheter such as a \"Foley\" catheter. Long-term use of indwelling urinary catheters in nursing home patients is well known as a potential cause of significant morbidity due to such infections.\nThis problem is widely recognized and many solutions for this problem have been suggested in the past. None of these solutions, however, have been completely free of secondary complications and/or completely successful in eliminating the problem. For instance, systemic use of antibacterial drugs or agents have been tried. However, these drugs generally have undesirable secondary effects upon the patient when used systemically, especially when there is a chronic need for catheterization and the drug must employed for a long period of time. Local use of such drugs or agents can be effective for a short period of time, but has not been found to be effective for long-term use for a number of reasons. First, the drug or agent is easily washed out if there is a leakage of urine through the urinary track outside of the catheter. Second, the drug or agent may be delivered only to certain areas of the urinary tract and, third, the drug or agent may be absorbed by the body tissues adjacent to the catheter and transported elsewhere within the body.\nOther methods of preventing periurethral contamination have been suggested. These include careful cleansing of the periurethral area on a routine basis, impregnating a sponge or the like with an antiseptic solution and retaining it in a position proximate the urethral, apply antimicrobial ointments to an external portion of the urinary tract, intraurethral irrigation of the urinary track through a separate channel, lubrication of the catheter with an antibiotic-containing material and the use of catheters impregnated with antimicrobial agents. Each of these methods has been explored and none have been found to be entirely satisfactory. In vitro tests of impregnated catheters indicate that the antibacterial agents within the catheters have a tendency to leach or diffuse out of the catheters in a short period of time. Often, the antibacterial activity was either gone or markedly diminished within 24 to 48 hours of insertion within the urethral tract. Therefore, it would be appreciated that a sustained release bactericidal cannula or catheter is needed in order to address the needs of patients requiring long-term catheterization or the like.\nAccordingly, it will be appreciated that there is a need for a medical device, cannula or catheter which will address these and other problems associated with the prior art devices. The present invention provides advantages over the prior art cannulas and catheters, over the prior art methods for manufacturing the same, and also offers other advantages over the prior art and solves other problems associated therewith."}
{"text": "As schemes that encode 3D video images, various schemes are proposed. In this case, the 3D video image is a video image generated by a video signal configured by a video signal of a first view point (to be referred to as a first-view-point video signal hereinafter) and a video signal of the second view point (to be referred to as a second-view-point video signal hereinafter). One of the first-view-point video signal and the second-view-point video signal is a video signal for a right eye, and the other is a video signal for a left eye. A video signal configured by only a video signal of a first view point is called a 2D video image.\nAs a scheme for encoding a 3D video image, for example, a scheme in which a first-view-point video signal is encoded with the same scheme as that for a 2D video image and a second-view-point video signal is subjected to motion compensation by using a picture of a first-view-point video signal simultaneous with the second-view-point video signal as a reference picture is proposed (to be referred to as a disparity compensation scheme hereinafter) is proposed. In comparison with the encoding for a 2D video image, the scheme can advantageously encode a 3D video image without decreasing the resolutions of a first-view-point video signal and a second-view-point video signal.\nAs another example of the scheme for encoding a 3D video image, a scheme that horizontally reduces a first-view-point video signal and a second-view-point video signal into halves, respectively, arranges the reduced video signals from side to side, and encodes the video signals by the same scheme as that for a 2D video image is proposed. In this case, information representing that an encoding stream is a 3D video image is added to the header information of the encoding stream to discriminate a 2D video image encoding stream from a 3D video image encoding stream.\nPatent Document 1 of a conventional technique describes an encoding apparatus that calculates quantized values such that a quantized value of an image (large parallax) on a front side is reduced and a quantized value of an image (small parallax) on a rear side is increased."}
{"text": "Adjustable hood-stops are assembly units disposed on the front-end bolster of vehicles. These units are often designed to prevent any impact between the vehicle hood and the vehicle headlights. Further, the use of adjustable hood-stops may enable the transfer of forces upon closure of the vehicle hood or upon vehicle impact. Depending on the specific vehicle model, the front-end bolster may contain at least two adjustable and two non-adjustable hood stops. These adjustable hood-stop units include a retaining portion, that may be integrated or detachable, on the body of the front-end bolster that receives a threaded bumper stop piece."}
{"text": "The present invention relates to a fungicidal mixture which comprises\na) a phenyl benzyl ether derivative of the formula I.a, I.b or I.c, \nxe2x80x83and/or a carbamate of the formula Id \nwhere X is CH or N, n is 0, 1 or 2 and R is halogen, C1-C4-alkyl or C1-C4-haloalkyl, it being possible for the radicals R to be different if n is 2, or a salt or adduct thereof, and\nb) a N-acetonylbenzamide of the formula II \nxe2x80x83where:\nR1 and R3 independently of one another are each halogen or C1-C4-alkyl;\nR2 is cyano, C1-C4-alkyl, C2-C4-alkenyl, C2-C4-alkynyl or C1-C4-alkoxy;\nR4 is hydrogen or C1-C4-alkyl;\nR5 is C2-C4-alkyl;\nR6 is thiocyano, isothiocyano or halogen, or a salt or adduct thereof,\nin a synergistically effective amount.\nThe invention further relates to methods for controlling harmful fungi using mixtures of the compounds I (I.a, I.b and I.c) and II and to the use of the compound I and the compound II for preparing such mixtures.\nThe compounds of the formula Ia-c, their preparation and their activity against harmful fungi are known from the literature (EP-A 253 213; EP-A 254 426; EP-A 398 692).\nThe compounds of the formula Id, their preparation and their activity against harmful fungi are known from the literature (WO-A 93/15,046; WO-A 96/01,256 and WO-A 96/01,258).\nAlso known are synergistic mixtures of the compounds II with dithiocarbamates, their preparation and their activity against harmful fungi (EP-A 753 258; U.S. Pat. No. 5,304,572).\nIt is an object of the present invention to provide mixtures which have an improved activity against harmful fungi combined with a reduced total amount of active ingredients applied (synergistic mixtures), with a view to reducing the application rates and to improving the activity spectrum of the known compounds I and II.\nWe have found that this object is achieved by the mixture defined at the outset. In addition, we have found that better control of harmful fungi is possible by applying the compound I and the compound II simultaneously, separately as well as together, or by applying the compound I and the compounds II in succession, than when the individual compounds are used.\nThe formula Id in particular represents carbamates in which the combination of the substituents corresponds to a line of the Table below:\nParticular preference is given to the compounds I.12, I.23, I.32 and I.38.\nThe formula II in particular represents N-acetonylbenzamides in which the combination of the substituents corresponds to a line of the Table below:\nParticular preference is given to those N-acetonylbenzamides which are generally and particularly preferred in EP-A 753 258.\nOwing to the basic character of their nitrogen atoms, the compounds Id and II are capable of forming adducts or salts with inorganic or organic acids or with metal ions.\nExamples of inorganic acids are hydrohalic acids such as hydrofluoric acid, hydrochloric acid, hydrobromic acid and hydriodic acid, sulfuric acid, phosphoric acid and nitric acid.\nSuitable organic acids are, for example, formic acid, carbonic acid and alkanoic acids, such as acetic acid, trifluoroacetic acid, trichloroacetic acid and propionic acid, and also glycolic acid, thiocyanic acid, lactic acid, succinic acid, citric acid, benzoic acid, cinnamic acid, oxalic acid, alkylsulfonic acids (sulfonic acids having straight-chain or branched alkyl radicals of 1 to 20 carbon atoms), arylsulfonic acids or aryldisulfonic acids (aromatic radicals, such as phenyl and naphthyl, which carry one or two sulfo groups), alkylphosphonic acids (phosphonic acids having straight-chain or branched alkyl radicals of 1 to 20 carbon atoms), arylphosphonic acids or aryldiphosphonic acids (aromatic radicals, such as phenyl and naphthyl, which carry one or two phosphoric acid radicals), it being possible for the alkyl or aryl radicals to carry further substituents, eg. p-toluenesulfonic acid, dodecylbenzenesulfonic acid, salicylic acid, p-aminosalicylic acid, 2-phenoxybenzoic acid and 2-acetoxybenzoic acid, etc.\nSuitable metal ions are, in particular, the ions of the elements of the first to eighth sub-group, in particular chromium, manganese, iron, cobalt, nickel, copper, zinc and of the second main group, in particular calcium and magnesium, and of the third and fourth main group, in particular aluminum, tin and lead. The metals can exist in the various valencies which they can assume.\nWhen preparing the mixtures, it is preferred to employ the pure active ingredients I and II, to which further active ingredients against harmful fungi or other pests, such as insects, arachnids or nematodes, or else herbicidal or glineth-regulating active ingredients or fertilizers can be admixed.\nThe mixtures of the compounds I and II, or the simultaneous joint or separate use of the compounds I and II, have outstanding action against a wide range of phytopathogenic fungi, in particular from the classes of the Ascomycetes, Basidiomycetes, Phycomycetes and Deuteromycetes. Some of them act systemically and can therefore also be employed as foliar and soil-acting fungicides.\nThey are specially important for controlling a large number of fungi in a variety of crop plants, such as cotton, vegetable species (eg. cucumbers, beans, tomatoes, potatoes and cucurbits), barley, grass, oats, bananas, coffee, maize, fruit species, rice, rye, soya, grapevine, wheat, ornamentals, sugar cane, and a variety of seeds.\nThey are particularly suitable for controlling the following phytopathogenic fungi: Erysiphe graminis (powdery mildew) in cereals, Erysiphe cichoracearum and Sphaerotheca fuliginea in cucurbits, Podosphaera leucotricha in apples, Uncinula necator in grapevines, Puccinia species in cereals, Rhizoctonia species in cotton, rice and lawns, Ustilago species in cereals and sugar cane, Venturia inaequalis (scab) in apples, Helminthosporium species in cereals, Septoria nodorum in wheat, Botrytis cinerea (gray mold) in strawberries, vegetables, ornamentals and grapevines, Cercospora arachidicola in groundnuts, Pseudocercosporella herpotrichoides in wheat and barley, Pyricularia oryzae in rice, Phytophthora infestans in potatoes and tomatoes, Plasmopara viticola in grapevines, Pseudoperonospora species in hops and cucumbers, Alternaria species in vegetables and fruit, Mycosphaerella species in bananas and Fusarium and Verticillium species.\nFurthermore, they can be used in the protection of materials (eg. in the protection of wood), for example against Paecilomyces variotii. \nThe compounds I and II can be applied simultaneously, that is either together or separately, or in succession, the sequence, in the case of separate application, generally not having any effect on the result of the control measures.\nThe compounds I and II are usually applied in a weight ratio of 10:1 to 0.01:1, preferably 5:1 to 0.05:1, in particular 1:1 to 0.05:1.\nDepending on the nature of the desired effect, the application rates of the mixtures according to the invention are, in particular in agricultural crops, from 0.01 to 8 kg/ha, preferably 0.1 to 5 kg/ha, in particular 0.5 to 3.0 kg/ha.\nThe application rates of the compounds I are from 0.01 to 2.5 kg/ha, preferably 0.05 to 2.5 kg/ha, in particular 0.05 to 1.0 kg/ha.\nCorrespondingly, in the case of the compounds II, the application rates are from 0.01 to 10 kg/ha, preferably 0.05 to 5 kg/ha, in particular 0.1 to 2.0 kg/ha.\nFor seed treatment, the application rates of the mixture are generally from 0.001 to 250 g/kg of seed, preferably 0.01 to 100 g/kg, in particular 0.01 to 50 g/kg.\nIf phytopathogenic harmful fungi are to be controlled, the separate or joint application of the compounds I and II or of the mixtures of the compounds I and II is effected by spraying or dusting the seeds, the plants or the soils before or after sowing of the plants, or before or after plant emergence.\nThe fungicidal synergistic mixtures according to the invention, or the compounds I and II, can be formulated for example in the form of ready-to-spray solutions, powders and suspensions or in the form of highly concentrated aqueous, oily or other suspensions, dispersions, emulsions, oil dispersions, pastes, dusts, materials for broadcasting or granules, and applied by spraying, atomizing, dusting, broadcasting or watering. The use form depends on the intended purpose; in any case, it should guarantee as fine and uniform as possible a distribution of the mixture according to the invention.\nThe formulations are prepared in a manner known per se, eg. by adding solvents and/or carriers. It is usual to admix inert additives, such as emulsifiers or dispersants, to the formulations.\nSuitable surfactants are the alkali metal salts, alkaline earth metal salts and ammonium salts of aromatic sulfonic acids, eg. ligno-, phenol-, naphthalene- and dibutylnaphthalenesulfonic acid, and of fatty acids, alkyl- and alkylarylsulfonates, alkyl, lauryl ether and fatty alcohol sulfates, and salts of sulfated hexa-, hepta- and octadecanols, or of fatty alcohol glycol ethers, condensates of sulfonated naphthalene and its derivatives with formaldehyde, condensates of naphthalene or of the naphthalenesulfonic acids with phenol and formaldehyde, polyoxyethylene octylphenol ether, ethoxylated isooctyl-, octyl- or nonylphenol, alkylphenol polyglycol ethers, tributylphenyl polyglycol ethers, alkylaryl polyether alcohols, isotridecyl alcohol, fatty alcohol/ethylene oxide condensates, ethoxylated castor oil, polyoxyethylene alkyl ethers or polyoxypropylene [lacuna], lauryl alcohol polyglycol ether acetate, sorbitol esters, lignosulfite waste liquors or methylcellulose.\nPowders, materials for broadcasting and dusts can be prepared by mixing or jointly grinding the compounds I or II or the mixture of the compounds I and II with a solid carrier.\nGranules (eg. coated granules, impregnated granules or homogeneous granules) are usually prepared by binding the active ingredient, or active ingredients, to a solid carrier.\nFillers or solid carriers are, for example, mineral earths, such as silica gel, silicas, silica gels, silicates, talc, kaolin, limestone, lime, chalk, bole, loess, clay, dolomite, diatomaceous earth, calcium sulfate, magnesium sulfate, magnesium oxide, ground synthetic materials, and fertilizers, such as ammonium sulfate, ammonium phosphate, ammonium nitrate, ureas, and products of vegetable origin, such as cereal meal, tree bark meal, wood meal and nutshell meal, cellulose powders or other solid carriers.\nThe formulations generally comprise 0.1 to 95% by weight, preferably 0.5 to 90% by weight, of one of the compounds I or II or of the mixture of the compounds I and II. The active ingredients are employed in a purity of from 90% to 100%, preferably 95% to 100% (according to NMR spectrum or HPLC).\nThe compounds I or II, or the mixtures, or the corresponding formulations, are applied by treating the harmful fungi or the plants, seeds, soils, areas, materials or spaces to be kept free from them with a fungicidally effective amount of the mixture, or of the compounds I and II in the case of separate application.\nApplication can be effected before or after infection by the harmful fungi."}
{"text": "A. Field of the Invention\nThe present invention relates generally to the field of solid body construction, and more particularly to the fabrication of three-dimensional objects using particle deposition.\nB. Description of the Related Art\nIn recent years, computer-aided design (CAD) technology has gained widespread acceptance in numerous engineering fields. Typical CAD systems permit the operator to design two- or three-dimensional constructs on a computer screen, then perform real-time rotation, translation and scaling operations with the on-screen image.\nMany CAD systems also allow designation of a coordinate system, allowing the operator to relate each point of the two-or three-dimensional design to a corresponding point on the selected coordinate system. This capability facilitates interface to automated manufacturing systems, which construct or assist with construction of the finished design based on its spatial coordinate representation.\nOne method of automated manufacture is disclosed in U.S. Pat. No. 4,665,492, in which particles are sequentially deposited on a substrate according to the three-dimensional coordinate representation of the solid design. Several particle compositions and methods of deposition are described, including ejection of ceramic material in a heated environment, ejection of plastic material in an ionized atmosphere, and the use of slurry material in a sub-freezing environment. All of these techniques appear limited in their ability to produce constructs that contain oblique surfaces.\nThis difficulty stems from the need to maintain subjacent support of particles during the deposition process. Without such support, the angle between adjacent oblique layers is limited by the adhesion characteristics of the deposited material, as well as the size of the particles. Because particle mass varies with the cube of the radius while surface area depends on the square of the radius, very small particles may be needed for particular materials to be deposited at a desired angle. Thus, designs containing outwardly displaced features may be impossible to fabricate with presently available techniques, depending on the degree of displacement, fabrication material and particle size.\nA second limitation of the method described in U.S. Pat. No. 4,655,492 relates to control of the ultimate spatial position of the deposited particle. The system described in this reference comprises two or more working heads to locate the point of deposition in space. This configuration must accommodate three degrees of freedom, resulting in the need for great control over movement of the working heads and precise monitoring of their positions."}
{"text": "Lubricant on rotating magnetic storage discs act as a protective coating for preventing corrosion of the magnetic layers of the disc. During the operation of a hard-disc, the air-bearing (that holds the magnetic elements) flies over the disc. During this process, it interacts with the lubricant coating on the disc. Some of the these interactions result in transfer of lubricant from the disc to the air bearing, where it can migrate to different areas of the head due to the joint effects of pressure and shear stresses. A similar mechanism is responsible for the accumulation on the head of organic contaminants, which originate mainly from the motor oil used in the spindle. The transferred lubricant and contaminants can accumulate in specific areas of the air bearing and eventually drop onto the disc. Such phenomena can affect reliability of the disc and in some cases even cause failures. Therefore, there remains a need to manage lubricant/contaminant on the air bearing."}
{"text": "The invention relates to a method for operating an internal combustion engine especially of a motor vehicle wherein the fuel, which is to be injected for a combustion, is directly injected into a combustion chamber during an intake phase and during a compression phase. Furthermore, the invention relates to an internal combustion engine especially for a motor vehicle having an injection valve with which the fuel, which is to be injected for a combustion, can be injected directly into a combustion chamber during an intake phase and during a compression phase. The engine also includes a control apparatus for the open-loop control and/or closed-loop control of the fuel mass injected into the combustion chamber.\nSystems for the direct injection of fuel into the combustion chamber of an internal combustion engine are generally known. In these systems, one differentiates a so-called stratified operation and a so-called homogeneous operation. The stratified operation is especially used for lower loads; whereas, the homogeneous operation is used for larger loads applied to the internal combustion engine.\nIn the stratified operation, the fuel is injected into the combustion chamber during the compression phase of the engine so that a fuel cloud is located in the immediate vicinity of a spark plug at the time point of the ignition. This injection can take place in different ways. Thus, it is possible that the injected fuel cloud is located at the spark plug already during the injection or directly after the injection and is ignited thereby. Likewise, it is possible that the injected fuel cloud is guided to the spark plug by a charge movement and is only then ignited. In both combustion processes, no uniform fuel distribution is present. Instead, there is a stratified charge.\nThe advantage of the stratified operation is that the applied lower loads can be handled by the engine with only a very small quantity of fuel. Larger loads can, however, not be satisfied by the stratified operation.\nThe fuel is injected during the intake phase of the engine in homogeneous operation provided for such larger loads so that a swirling and therefore a distribution of the fuel can still easily take place in the combustion chamber. To this extent, the homogeneous operation approximately corresponds to the operation of engines wherein fuel is injected into the intake manifold in the conventional manner. The homogeneous operation can be utilized, as required, also for lesser loads.\nIn the stratified operation, the throttle flap is opened wide in the intake manifold leading to the combustion chamber and the combustion is open-loop controlled and/or closed-loop controlled essentially only by the fuel mass to be injected. In the homogeneous operation, the throttle flap is opened or closed in dependence upon the requested torque and the fuel mass, which is to be injected, is open-loop controlled and/or closed-loop controlled in dependence upon the inducted air mass.\nIn both modes of operation, that is, in the stratified operation and in the homogeneous operation, the fuel mass, which is to be injected, is open-loop controlled and/or closed-loop controlled additionally in dependence upon a plurality of additional input quantities to an optimal value with respect to fuel saving, exhaust-gas reduction and the like. The open-loop control and/or closed-loop control is then different in the two modes of operation.\nIt is possible to combine the two described modes of operation with each other. This is known as a so-called double injection. For this purpose, the fuel mass is determined which in total, is to be injected. Then, a first part of the fuel mass to be injected is injected during the intake phase and the remaining second part of the fuel mass to be injected is injected into the combustion chamber directly thereafter during the compression phase. The fuel mass, which is to be injected, for one and the same combustion is therefore injected into the combustion chamber of the engine with the aid of two injections.\nThe advantage of a double injection of this kind is that an abrupt switchover between the two modes of operation is avoided, that is, a switchover between the homogeneous operation and the stratified operation and back. In lieu thereof, the switchover is xe2x80x9cblurredxe2x80x9d over a wide range in that, for example, the fuel mass, which is to be injected during the intake phase, is slowly reduced and the fuel mass which is to be injected during the compression phase, is slowly increased.\nLikewise, it is possible to provide a tank venting for the engine. Here, the so-called regeneration gas, which is vaporized in a fuel tank, is collected and supplied to the combustion. The advantage of a tank venting of this kind is that the fuel, which is contained in the regeneration gas, does not escape to the ambient but is used by the engine.\nThe object of the invention is to provide a method for operating an internal combustion engine as well as to provide an internal combustion engine wherein the advantages of the double injection and the tank venting are present.\nThis object is achieved in accordance with the invention in a method of the kind referred to above in that a regeneration gas is supplied to the combustion chamber during the intake phase. The object according to the invention is achieved in an internal combustion engine of the kind initially mentioned above in that a regeneration gas can be supplied to the combustion chamber during the intake phase by the control apparatus.\nBy supplying the regeneration gas during the intake phase, the breakdown of the entire quantity of fuel to be injected to the intake phase and the compression phase is not affected or even disturbed. The accuracy especially of the fuel mass, which is to be injected during the compression phase, is therefore maintained undiminished. The advantage is achieved notwithstanding that the regeneration gas is supplied to the combustion chamber and therefore to the combustion. The invention therefore combines the advantages of the direct injection (especially of the reduced fuel consumption and the lower discharge of toxic substances of the stratified operation) with the advantages of the tank venting, especially the complete utilization also of vaporizing fuel.\nIt is especially advantageous when fuel is injected and regeneration gas is supplied to the combustion chamber during the intake phase and when fuel is injected into the combustion chamber during the compression phase but no regeneration gas is supplied.\nIn an advantageous embodiment of the invention, the fuel mass, which is supplied in total into the combustion chamber during a single combustion, is composed of a first part and a second part. The first part is supplied during the intake phase and the second part is injected during the compression phase. The first part is put together from a third injected part and from a fourth part comprising the regeneration gas.\nIn an advantageous embodiment of the invention, for the first part of the total mass of supplied fuel, the third injected part is reduced and the fourth part made up of the regeneration gas is increased.\nIn this way, it is achieved that as much vaporized fuel as possible is supplied to the combustion chamber of the engine. The use of the vaporized fuel is thereby optimized.\nIt is especially advantageous when so much regeneration gas is supplied to the combustion chamber that the third injected part reaches a minimum value or zero.\nThe realization of the method of the invention in the form of a control element is especially significant with this control element being provided for a control apparatus of an internal combustion engine, especially of a motor vehicle. A program is stored on the control element which is capable of being run in a computer apparatus and especially in a microprocessor and is suitable for carrying out the method of the invention. In this case, the invention is realized by a program stored on the control element so that this control element, which is provided with the program, can define the invention in the same manner as the method for whose execution the program is suitable. As a control element, especially an electrical storage medium can be applied, for example, a read-only-memory.\nAdditional features, application possibilities and advantages of the invention become apparent from the following description of the embodiments of the invention which are shown in the figures of the drawing. Here, all described or illustrated features constitute the subject matter of the invention by themselves or in any desired combination independently of their combination in the patent claims or their dependency as well as independent of their formulations or illustration in the description or in the drawing."}
{"text": "This invention relates to a process for improving the ozone resistance of certain vulcanized polymers and polymer blends, and in some cases also their flex resistance, as well as to polymers and polymer blends having improved ozone resistance and sometimes also flex resistance.\nElastomeric chloroprene polymer vulcanizates are quite resistant to ozone at normal temperature but are more readily degraded at elevated temperatures, especially when flexed or stretched. Other polymer vulcunizates are less resistant to ozone than chloroprene polymer vulcanizates. Various antiozonants are commercially available and usually are incorporated in such polymer formulations. While many such additives provide a good degree of protection against ozone deterioration, certain commercial antiozonants cause premature curing of polymer stocks, while most cause discoloration or staining of the polymer vulcanizates. There presently appears to be only one nonstaining antiozonant commercially available.\nU.S. Pat. No. 3,563,947 (to W. F. Gruber) discloses the use of certain norbornene derivates as antiozonants for elastomeric chloroprene polymers. The copending allowed application Ser. No. 534,200 of M. C. Chen describes certain triarylphosphines which improve the flex resistance of chloroprene-sulfur copolymers."}
{"text": "1. Field of the Invention\nThe present invention is directed to a method for data entry into a postage meter machine, and to an apparatus for franking postal matter and for producing a franking advert mark uniquely allocated to a data center which communicates with the meter to alter the contents of accounting registers in the meter.\n2. Description of the Prior Art\nA postage meter machine is utilized for franking postal matter and is equipped with at least one input means, particularly with a chip card write/read unit, an input/output control means and with an output means.\nChip cards are known which comprise a plurality of non-volatile memories, or separately accessible memory areas and a microprocessor in order to transmit data representing different types of information into the postage meter machine and in order to read data out of the postage meter machine.\nA postal fee accounting system disclosed in U.S. Pat. No. 5,111,030, corresponding to German OS 39 03 718 includes a write/read unit for chip cards that is connected to a postage meter machine via a control unit. Data about postage meter machine use are written into the chip card, this data being read out later with a personal computer for accounting.\nGerman OS 40 33 164 discloses a metered postage tape system having chip card write/read unit in a common housing. Two memory areas exist in the non-volatile data memory of the chip card, the first memory area containing the current postage credit data and the second memory area containing the data for an individual advertising design (advert mark). The chip card assigned to a user can also be used in a known way for loading or recharging the postage meter machine and has a transport and accounting function for the exchange of data for fee accounting as well as a reloading function for an advertising design. The second memory area for the individual advertising design can be selected in the chip card separately from the first memory area for the postage credit and is only read when needed.\nIt would be desirable to additionally store accounting data (date, number and value of impressions) in the chip card in order to be able to undertake a chronologically following accounting with a personal computer. However, the storage space of this known chip card is completely occupied by the additional image data, so that no further data can be stored.\nThis solution is also time-consuming because, after the entry of an identification code for each user, an individual advertising design must be newly loaded from the chip card into the processor system (postage meter) via the serial interface and the old, stored, individual advertising design in the processor system must be replaced by a new, individual advertising design. The data exchange procedure via a serial interface already lasts longer than the data exchange via a parallel interface. The advertising design can thus not be constantly reloaded for every impression. This solution is thus ineffective given a constant or frequent change of users.\nDue to the limited storage space that is available on a chip card, a plurality of chip cards simultaneously plugged into a plurality of write/read units would again have to be employed for further data to be entered, this further increasing the time consumed in a data read-out associated with every printing event.\nU.S. Pat. No. 4,812,994 discloses a system intended to prevent unauthorized access to use the postage meter machine by inhibiting the postage meter machine given the absence of an identification signal ID and/or after the passage of a predetermined time interval without use. The ID signal can be entered by a chip card, by a personal computer, via modems or can be manually entered into the postage meter machine. The postage meter machine is enabled after a positive comparison with a user identification signal stored in the postage meter machine.\nIn the aforementioned solutions, the postage meter machines are inhibited for use until the chip card is plugged into a corresponding write/read unit, as a result of which the authorization of the user is initiated.\nWO93/05482 corresponding to German OS 41 29 302 proposes a modified solution for incrementing the fee credit in the credit balance memory of postage meter machines on the basis of a chip card that carries a reloading credit that, when subsequently erased, can in turn be removed. In another version, the credit stored in the chip card is debited step-by-step. However, no further data can be fetched from the chip card. A pluggable EPROM is introduced into a permanently installed plug-in socket for the postal fees (postage fee table).\nIt is standard for service to install the new postage fee values at the customer on the basis of a non-volatile memory module. In addition to the substantial service outlay for this procedure, this may result in the use of an out-of-date postage fee table because of the necessity to replace the table in advance of its effective date, or after its effective date, since, with a large number of meters, every one cannot be installed with a new module on the same day.\nU.S. Pat. No. 3,635,297 discloses a mail-handling apparatus having automatic fee value calculation. A replaceable memory (ROM) contains a postage fee table from which the amount of postage is calculated upon entry of the weight of an item to be mailed and possibly with the entry of further mailing particulars. A disadvantage of this approach is the outlay that arises when, due to a new fee schedule taking effect, the fee values must be modified in the table of the memory. It is not assured that the user of the apparatus will have the replacement of the memory for updating undertaken in time.\nWhen a new schedule of fees takes effect, U.S. Pat. No. 4,122,532, corresponding to German OS 28 03 982 proposes a remote valuation for updating the postage fee table stored in postage meter machines. Under remote actuation, the new fee schedule is simultaneously transmitted to a plurality of postage meter machines from a central data station. Such an updating with the new fee schedule, however, assumes that the postage meter machines are turned on at these points in time and can be constantly addressed.\nFor protection against fraudulent manipulations, U.S. Pat. No. 4,933,849, corresponding to German OS 38 23 719 discloses that a representative character pattern be printed out beginning with a specific date. When examining the mail, the printed date and the character are compared in the Post Office to the pattern that is authorized for this date. An authorization means which comprises a memory means for storing data representing a plurality of character patterns and dates serves the purpose of printing. The data that allocate the representative character pattern to a defined date are updated via a remote valuation with an external selection means when the users of the postage meter machines request a recharging. This security system, however, is restricted to point-to-point networks and cannot be applied to portable postage meter machines that are carried from one location to another (mobile office).\nIt has thus not been previously possible to update portable postage meter machines, i.e. postage meter machines that are not permanently installed via a telephone network and to secure these against fraudulent manipulations.\nU.S. Pat. No. 4,506,330 discloses a removable printed circuit board having DIP switches at the address input of a PROM for the selective addressing of a memory sector. Mail differences within special service classes can be economically calculated. Although the storage space for the customer data storage of postal zone tables is thereby significantly reduced, the DIP switches must be manually set.\nBy contrast, the postage meter machine disclosed in U.S. Pat. No. 4,138,735 comprises a rate PROM that can be equipped with a current postal fee table by radio or telephone. The possibility of a non-simultaneous transmission of a postage fee table is achieved on the basis of an individual addressing of each postage meter machine. However, it is complicated to assure that the contacting on the part of the central data center was successful. Another disadvantage is that it is unavoidable that the memories of the postage meter machine are filled with many unnecessary data.\nAs much data as would be desirable can definitely not be currently stored in a single chip card. Given an increased data processing outlay, however, it would be possible to reduce the total data to a \"necessary\" data set and to store only the necessary data set in a chip card.\nAnother way of bypassing the limited storage capacity present on a chip card would be to employ a plurality of chip cards simultaneously plugged into write/read units.\nU.S. Pat. No. 4,802,218 discloses an automatic transmission system having a plurality of slots for chip cards that, in addition to employing a chip card for recharging credits and for accounting, whereby the postage fee value to be printed is subtracted from the credit, also simultaneously employs a further chip card for a postage fee table with whose assistance the aforementioned postage fee value is calculated. As a result of the plurality of write/read units, however, the apparatus becomes too large and too expensive.\nOn the other hand, accounting (debiting) data are already stored in the postage meter machine, but heretofore must additionally be stored in a chip card in order to transport the accounting data to a personal computer equipped with a chip card reader and to print out an accounting log via a printer connected thereto. However, variable printing formats can be produced with an electronic printer in postage meter machines.\nIt is disclosed in pending German Patent Application P 42 24 955.4 also to realize the printing of lists, for example internal accounting reports about the use of the credit stored in the postage meter machine from the individual cost allocation accounts (a customer may have a plurality of cost allocation accounts), on the basis of electronic printing processes, for example on the basis of a thermal transfer printing process. With these enhanced possibilities, however, the operation of the postage meter machine via the keyboard can easily become so complicated that an unqualified user cannot not undertake it.\nConventional automatic franking machines have input possibilities in order to set values or prescriptions or in order to input commands. Many keys are either used or few keys must be multiply occupied and be sequentially actuated.\nIn a known postage meter machine commercially available from Francotyp Postalia GmbH, a number is allocated to every advertising design electronically stored in the machine. After the selected number has been fetched by pressing a key, a function key for setting the advert mark is actuated in order to modify the advertising design in accord with the selected number.\nSince it is preferable that the outlay for the operating elements should not increase, an inexpensive, faster possibility of operation is required given a simple user interface. A non-authorized user of the postage meter machine should be prevented from fetching the data of other cost allocation accounts merely by pressing a button."}
{"text": "The present invention relates to a device to be used as a connector or a repair device for pipe or tubing. More particularly, the present invention relates to a device which utilizes a resilient sheet member which applies an inwardly contracting force on the exterior surface of the pipe or tubing.\nThere has been a need in the past for a cost effective, efficient and rapid means of connecting together pipe or tubing and for the repairing of faults or holes in such pipe or tubing without the necessity of cutting or dividing the pipe or tubing. There has been a need in the past for a connector which does not require the insertion of an inner supporting member within the pipe or tubing to be connected together or repaired. Furthermore, there has been a desire to provide a strong adhesive bond between the connector and the pipe or tubing to be bonded together or repaired.\nThe present invention provides advantages directed to providing solutions to these needs."}
{"text": "1. The Field of the Invention\nThe present invention relates to the field of telephone network and device testing. Specifically, the present invention relates to using a call control server to test a telephone network or device in an automated and integrated fashion.\n2. Background and Relevant Art\nThe telephone has been one of the most pervasive inventions of the modern era. It is now common practice for an individual to use a telephone to audibly converse in real-time with another individual even if the caller and the callee are remotely located. Recently, the types of devices able to support telephonic communication have expanded well beyond what was conventionally thought of as a telephone. For example, digital telephones (hardwired and wireless), and some types of personal computers and Personal Digital Assistants are now able to engage in telephonic communication.\nConventionally, telephonic device manufacturers typically test their telephonic devices prior to distribution to be sure the telephonic device works as intended. In order to ensure the telephonic device works during normal use, it is desirable to test the telephonic device under a wide variety of possible conditions that the telephonic device is anticipated to experience during normal operation. However, normal operation of a telephonic device involves connection and communication with other telephonic devices. Thus, the telephonic device cannot be tested in isolation.\nOne way to test the operation of a telephonic device and associated telephonic network is to actually have the subject telephonic device communicate in various scenarios with other surrounding telephonic devices, and evaluate the performance of the subject telephonic device and network under the testing scenarios. While the variety of such scenarios are innumerable, such scenarios include, for example, receiving a ring generated by another telephonic device calling in, answering the telephone call and communicating between the devices, engaging in a conference call, or similar common telephonic activities.\nAccording to one conventional method, an individual stations the other surrounding telephonic devices in order to provide the necessary input to those other telephonic devices in order to test the subject telephonic device. However, this requires that a human user take the effort to enter in the correct sequences of input at the correct times in order to provide the desired testing conditions for the subject telephonic device.\nIn order to remove the cost and unpredictability associated with using humans to station the surrounding telephonic devices, server technology has developed that allows a server to emulate the associated telephone network and surrounding telephonic devices. Servers that have this capability are often called “call control” servers.\nConventional call control servers allow a tester to dial into the call control server and instruct the call control server to perform certain telephone related actions. The instructions may be dispatched by first establishing a telephonic connection with one of the lines of the call control server. Then, the tester may dispatch the instructions by dialing the keypad of a conventional telephone with a particular audible tone sequence. Specifically, when the user enters a digit such as 0, 1, 2, 3, 4, 5, 6, 7, 8, 9, * or # on the keypad, a distinct audible tone is asserted on the telephone connection. The call control server converts the frequency of each successive tone into its corresponding digit 0 through 9, * or #. Specifically, the tones are generated and interpreted using a known standard such as the Dual Tone Multi-Frequency. Then, the call control server may interpret the digits and implement the associated instruction. Thus, a tester may control a conventional call control servers via standard telephone voice lines.\nWhile conventional call control servers are useful in eliminating the need for a person to station surrounding telephones during testing, the number of commands that may be implemented is fairly limited. This is due, in part, to the limited way in which a tester issues commands; namely, through DTMF tones. The dial tone command set associated with a particular call control server is typically not intuitive since the command set is limited to the digits 0 through 9, * and #. Since the command set is not intuitive, it is difficult for a human tester to remember a large number of commands. Accordingly, the amount of scenarios enabled by conventional call control servers is also relatively limited.\nMethods also exist for issuing non-audible data commands to call control servers over a data line. However, these conventional methods are limited in the level of integration between the issuing of data commands over data lines, and the actual implementation of such commands on voice lines. For example, under conventional technology, one cannot issue a data command over a data line to command a call control server to establish a multi-line conference call over voice lines.\nTherefore, what is desired are mechanisms for telephonic testing in which the call control server has more integrated data and voice features, and in which the tester may use a more intuitive command set thereby allowing for more testing scenarios."}
{"text": "1. Field of the Invention\nThe present invention relates to an exposure apparatus.\n2. Description of the Related Art\nLithography is a well-known technique for forming a predetermined circuit pattern on a semiconductor. Lithography includes a step of forming a latent image pattern by irradiating a substrate coated with a photosensitive agent (photoresist) with light via a reticle to expose the photosensitive agent, and, thereafter, developing the photosensitive agent. Development forms a mask pattern for use in, e.g., etching.\nAs indicated by equation (1), the resolution of the exposure apparatus is known to be proportional to a wavelength λ of a light source and inversely proportional to a numerical aperture NA of a projection lens.Resolution=k1·(λ/NA)  (1)where k1 is a proportionality constant.\nTo improve the resolution of the exposure apparatus, it is only necessary to shorten the wavelength of the light source or to increase the numerical aperture of the projection lens.\nOne of the characteristic features of an optical system of the exposure apparatus is the depth of focus (DOF). The DOF is the allowable blurring range of a projected image, which is expressed by the distance from the focus point. The DOF is given by:DOF=k2·(λ/NA2)  (2)where k2 is a proportionality constant.\nAs the wavelength of the light source shortens or the numerical aperture of the projection lens increases, to improve the resolution of the exposure apparatus, the DOF noticeably decreases. This shortens a distance, in the optical axis direction, in which accurate processing is possible.\nEspecially, next-generation devices aiming at an increase in the degree of integration by forming a finer, stereoscopic circuit pattern are facing a serious problem of such a decrease in the DOF. This is because a wide range of sharp focus is necessary, since the formation of a stereoscopic circuit pattern elongates the processing dimension in the optical axis direction and a given DOF is always necessary irrespective of the degree of micropatterning of the circuit.\nTo solve the above problem, an attempt to increase the DOF is made by projecting the pattern of the reticle onto the substrate using exposure light, including a plurality of wavelengths, to form images at different positions on the same optical axis.\nFor example, Japanese Patent No. 02-619473 proposes an exposure method of using, as an exposure light beam, the light obtained by synthesizing the light beams from both a light source, which emits the first wavelength light and a light source, which emits the second wavelength light.\nJapanese Patent Laid-Open No. 11-162824 proposes a method of performing exposure with exposure light, including a plurality of wavelengths, by inserting, on the light path between the light source and the wafer, a filter, which selectively transmits light in a plurality of wavelength ranges.\nThe current exposure apparatus must stabilize the light quality of the light source at the time of exposure and its calibration.\nFIG. 1 shows the spectral shape of light (single-wavelength light) emitted by the light source. The abscissa represents a light wavelength λ, and the ordinate represents the light intensity as a function of the wavelength λ.\nAs shown in FIG. 1, the current exposure apparatus employs a central wavelength λ0, FWHM (Full Width Half Maximum), and an E95 value as evaluation indices of the light quality. The FWHM indicates a bandwidth measured at an intensity at half of the peak intensity. The E95 value indicates a bandwidth in which 95% of the entire spectral energy is concentrated.\nThe exposure apparatus monitors the central wavelength λ0, FWHM, and E95 value immediately before exposure or during exposure to calculate a variation (3σ) in measurement values and an error value (an error with respect to a command value), thereby confirming the stability of the light quality.\nEven exposure using exposure light including a plurality of wavelengths requires stabilization of the light quality at the time of exposure and calibration of the exposure apparatus.\nHowever, light including a plurality of wavelengths exhibits a spectrum, such as that shown in FIG. 2. Even when the central wavelength λ0, FWHM, and E95 value as conventional indices are calculated for the entire spectrum of light including a plurality of wavelengths, it is impossible to sufficiently confirm the light quality based on these indices.\nFor example, when the pattern of the reticle is projected onto the wafer with light including two peaks, as shown in FIG. 2, an adverse effect, such as image blurring, occurs. To prevent this problem, a difference Ediff between the peak intensities should also be controlled/monitored during exposure."}
{"text": "This invention relates to a code break mechanism for a stacking apparatus receiving a stream of articles, more particularly to a code break mechanism for producing discontinuities in the stacker article stream in accordance with a selected parameter changes in the article stream.\nIn article processing systems, generally the last process step is the collection of processed articles. For example, in a conventional mail processing system, a power stacker may be employed to collect envelopes ejected from a mailing machine. Generally, the power stacker receives the envelopes in a seriatim and uniform manner and stacks the envelopes. An operator periodically removes the accumulated envelope stack from the stacker for packaging in preparation for subsequent deposit with a mail forwarding agent, such as, the United States Postal Service (USPS).\nMail forwarding agents, such as, the USPS, offer to the mailers special postage discount rates for mail which is presorted, for example, by zip code. Therefore, where the mail is processed in a presorted manner, for example, according to zip code, it would be advantageous to the operator for the power stacker to include a means for visually marking zip code changes or breaks within the stacker received mail stream."}
{"text": "1. Field of the Invention\nThe present invention relates to flexible magnetic sheets, which are applied to the skin of persons or animals for therapeutic or analgesic purposes.\n2. Background Of The Invention\nMagnetic therapeutic plasters are described in U.S. Pat. No. 4,489,711 to Latzke (\"Latzke\"), which is incorporated by reference herein. The magnetic poles in FIG. 1 of Latzke are disposed in alternating rows of \"+\" and \"-\" (i.e., North and South) poles. Latzke implemented the invention on magnetic rubber sheets having a thickness of 0.5 to 1.5 mm thick, with the magnetic poles 5-10 mm apart, and a magnetic field strength of 50 to 10,000 gauss, preferably, 400 to 2,000 gauss. The examples in Latzke describe the beneficial effects of treatment with the magnetic plasters.\nU.S. Pat. No. 4,549,532 to Baermann (\"Baermann\") discloses several different geometries for a flexible magnetic sheet, including circular strips of alternating polarities (FIG. 1), triangular strips (FIG. 2), rectangular or square strips (FIG. 3) and octagonal strips (FIG. 4). U.S. Pat. Nos. 5,277,692, 5,538,495 and 5,514,072 to Ardizzone disclose a flexible magnetic pad in which alternating N/S magnets spiral out from a central core such that triangular wedges of N/S alternating magnets radiate from the central core. U.S. Pat. No. 5,304,111 to Mitsuno et al. disclose a flexible magnetic sheet in which the N/S magnetic poles are disposed in a checkerboard pattern. Other U.S. patents of interest include U.S. Pat. Nos. 4,162,672 to Yazaki, U.S. Pat. No. 5,336,498 to Snider, U.S. Pat. No. 5,017,185 to Baermann and U.S. Pat. No. 5,259,892 to Kubota. Foreign patents of interest include French Patent No. 1,215,110 to Tanaka, Japanese Laid-Open Patent Application No. 56-7405 to Miyake, and French Patent 2,371,916 to Van Den Bulke."}
{"text": "As a method for imparting a plurality of functions to an acrylic polymer, a method is known of emulsion-polymerizing a plurality of layers having different functions in multiple steps.\nFor example, as an agent for modifying the impact-resistant of an epoxy resin, there is a report on a core-shell type powdery polymer having an inner layer composed of a (meth)acrylate polymer having a glass transition temperature of −30° C. or lower and an outer layer composed of a (meth)acrylate polymer having a glass transition temperature of 70° C. or higher (patent document 1).\nFurther, as a vinyl chloride sol alternate material, (meth)acrylic polymer particles are known composed of a core layer giving the flexibility of a coated film and a shell layer suppressing contact of the polymer of the core layer to a plasticizer, thereby giving storage stability. For example, there is a report on a core-shell type acrylic polymer particle composed of a core part and a shell part each having a specific solubility parameter and a specific glass transition temperature (patent document 2). Moreover, there is a report on a method of adding a polymerization inhibitor in the course of a polymerization step to produce an emulsion containing polymers having different particle sizes (patent document 3).\nIn the method described in patent document 1, however, if the mass ratio of the shell component in the polymer particle is reduced, it becomes difficult to form a particle covered uniformly by the shell layer. It is believed that this difficulty is caused since a polymerization reaction progresses sequentially from a part of contact of the monomer of the shell component to the polymer of the core. When the amount of the core component is 80% by mass or more, it becomes difficult to coat the core component completely by the shell layer, and storage stability thereof tends to be insufficient. If the shell component is used in an amount necessary for uniform coating of the core component, the proportion of the core component decreases, leading to a polymer in which the effect of modifying impact resistance as the original object is low.\nAlso in the method described in patent document 2, if the mass ratio of the shell part in the polymer particle is reduced, it becomes difficult to form a sufficient coat by the shell component. Thus, if the mass ratio of the core component is increased to obtain a coated film giving satisfactory flexibility, it becomes difficult to completely coat the core, and storage stability thereof tends to be insufficient. Moreover, if the shell component is used in an amount necessary for uniform coating of the core component, the flexibility of the coated film lowers.\nThe method described in patent document 3 is a method for obtaining a dispersion containing two polymer particles having different physical properties, and is not a method in which a polymer of a core component is coated uniformly by a small amount of a polymer of a shell component and the performances of the core component and the shell component are allowed to manifest.\nIn the conventional methods, as described above, it is difficult to uniformly coat a polymer of a core component by a small amount of a polymer of a shell component and to manifest the performances of the core component and the shell component in a good balance."}
{"text": "In the art of thermal ink jet printing, it is known to provide controlled and localized heat transfer to a defined volume of ink which is located adjacent to an ink jet orifice. This heat transfer is sufficient to vaporize the ink in such volume and cause it to expand, thereby ejecting ink from the orifice during the printing of characters on a print medium. The above predefined volume of ink is customarily provided in a so-called barrier layer which is constructed to have a plurality of ink reservoirs therein. These reservoirs are located between a corresponding plurality of heater resistor elements and a corresponding plurality of orifice segments for ejecting ink therefrom.\nOne purpose of these reservoirs is to contain the expanding ink bubble and pressure wave and make ink ejection more efficient. Additionally, the reservoir wall is used to slow down cavitation produced by the collapsing ink bubble. For a further discussion of this pressure wave phenomena, reference may be made to a book by F. G. Hammitt entitled Cavitation and Multiphase Flow Phenomena, McGraw-Hill 1980, page 167 et seq, incorporated herein by reference.\nThe useful life of these prior art ink jet print head assemblies has been limited by the cavitation-produced wear from the pressure wave created in the assembly when an ink bubble collapses upon ejection from an orifice. This pressure wave produces a significant and repeated force at the individual heater resistor elements and thus produces wear and ultimate failure of one or more of these resistor elements after a repeated number of ink jet operations. In addition to the above problem of resistor wear and failure, prior art ink jet head assemblies of the above type have been constructed using polymer materials, such as those known in the art by the trade names RISTON and VACREL.\nWhereas these polymer materials have proven satisfactory in many respects, they have on occasion exhibited unacceptably high failure rates when subjected to substantial wear produced by pressure waves from the collapsing ink bubbles during ink jet printing operations. Additionally, in some printing applications wherein the printer is exposed to extreme environments and/or wear, these polymer materials have been known to swell and lift from the underlying substrate support and thereby render the print head assembly inoperative."}
{"text": "Medical ultrasound imaging systems are using increasingly complex transducer elements to achieve higher quality images and measurements. The most sophisticated medical ultrasound scanners now use linear arrays (N.times.1) containing over one hundred transducer elements and two dimensional arrays (N.times.M) containing upwards of five hundred channels. The two-dimensional arrays with their large number of transducer elements offer several advantages :over the linear arrays for B-scans including; focusing in the elevation dimension, phase aberration corrections in two-dimensions, concurrent orthogonal B-scans, real time volumetric imaging, high speed C-scans and angle independent blood flow imaging.\nHowever, while the increased number of transducer elements have improved the imaging capability of the ultrasound system, the electronic cabling requirements have severely impacted the size, cost and flexibility of the transducer cables. Each transducer element, located in a probe which touches the patient body, requires an independent exciter circuit to supply an electronic shock-excitation or burst pulse input to generate an ultrasonic energy pulse. Currently, coaxial cables carry the electronic pulse from the ultrasound scanner to the independent transducer elements located in the transducer probe. The large number of independent transducer elements typically require a like number of independent coaxial cables. The coaxial cable with insulation and shield typically measures 0.46 mm diameter or greater per channel. The large quantity of channels multiplies the size of the transducer assembly cable which can make it extremely stiff, heavy, and costly.\nPliant fiber optic cables have diameters nearly ten times smaller than the coaxial cables. Fiber optic cables have not been used to transmit high voltage to transducers with a multiplicity of elements but have seen limited use in intra-cavity and intra-blood-vessel ultrasound imaging to electrically isolate the patient."}
{"text": "In many conveying devices such as, for example, elevator installations, cranes or hoists use is made of belt-shaped support means. These support means generally comprise a plurality of tensile carriers which consist of steel wires and which accept the tensile forces to be absorbed by the support means. The tensile carriers are generally surrounded by a casing of synthetic material. The casing protects the tensile carriers from, for example, mechanical wear, since the support means are frequently guided around deflection points. In addition, the casing improves the traction of the support means on deflecting or drive rollers and fixes the arrangement of the tensile carriers relative to one another.\nSuch support means are a safety-relevant component within a conveying device. Failure or breakage thereof can lead to dropping down of the items being conveyed. This can lead to substantial harm to objects or persons. For this reason, use is made of check units in conveying devices, which check, in particular, the mechanical state of the tensile carriers. Damage of the tensile carriers accepting the forces shall thereby be able to be recognized in good time so that the support means in the case of damage can be exchanged in order to prevent failure of the conveying device.\nThe electrically conductive, metallic tensile carriers are surrounded by the electrically insulating casing of synthetic material. In order to carry out a check of the state of the tensile carriers contacting of a contact element with the tensile carriers is necessary in some methods. In one known method an electrical current, which serves as a test current for the purpose of ascertaining the state of the tensile carriers, is conducted through the tensile carriers with the help of the contact element.\nDE 3 934 654 A1 shows a support means of the category described above. The ends of the tensile carriers are in that case conductively connected in pairs with a bridge member so that the tensile carriers of the support means are electrically connected in series. The tensile carriers of the support means are connected with a voltage source by way of an ammeter so that the state of the tensile carriers can be assessed by means of the test current which is conducted through all tensile carriers by virtue of the electrical connection in series.\nWO 2005/094249 A2 shows a system for contacting a support means in which the contact elements puncture the casing of the support means perpendicularly to a longitudinal axis of the tensile carriers and then penetrate into the tensile carriers. It can be disadvantageous in that case that the contact elements due to the required puncturing process through the casing miss the tensile carriers. In addition, a transfer resistance between the tensile carriers and the contact elements penetrated therein can change over time which has a disadvantageous effect on the meaningfulness of a monitoring procedure.\nWO 2010/057797 A1 and WO 2011/003791 A1 show systems for contacting a support means in which contact elements are mounted on exposed tensile carriers of the support means, for example by resilient contacts or by contact spikes penetrating into the tensile carriers. It is disadvantageous with such contacting systems that the connection between contact elements and tensile carriers is non-constant and that a high transfer resistance between the contact element and the tensile carriers is present."}
{"text": "The present invention relates to d.c. - a.c. inverter circuitry connectible at its input to a d.c. power source and at its output to a fluorescent lamp. More particularly, it relates to such circuitry that enables a standard type of fluorescent lamp to be operated from a low voltage battery with a high degree of efficiency, i.e., with minimum electrical power input for a given level of light power output.\nThere are many applications where it is desirable to have a light source that is capable of operating from a convenient d.c. power source, such as a battery. The incandescent lamp is often used as a means of converting electrical power from a battery into light power (illumination). It is well known in the art, however, that an incandescent lamp is a very inefficient means of converting electrical energy into visible light. Approximately 6% of the electrical energy supplied to an incandescent lamp is converted into visible light, the other 94% being used up in the generation of heat with a small accompanying amount of invisible ultraviolet light.\nWhen the supply of electrical energy is limited, as in the case of a battery-operated lamp, it is desirable to achieve maximum operating efficiency so that the lamp can be operated for the longest period of time before the supply of electrical energy is exhausted. It is also well known in the art that a fluorescent lamp is approximately four to five times more efficient than an incandescent lamp. Because of its higher efficiency, the fluorescent lamp has the potential to out-perform the incandescent lamp in those instances where efficiency ranks as an important operating parameter. In order to take full advantage of this remarkable potential, however, it is essential that transistor circuitry used to supply power to the fluorescent lamp does not itself operate inefficiently. To this end, it has heretofore been recognized that the circuitry should at least take into account that the gas mixture in a fluorescent lamp must first become highly ionized in order to cause the lamp to emit visible light, that the application across the lamp of an a.c. voltage rather than a d.c. voltage makes it possible to use a purely reactive element as a current limiting impedance (ballast) in order to minimize energy loss relative to the loss experienced with the use of a resistive ballast, and that the lamp impedance will remain fairly constant if the frequency of the a.c. voltage applied across the lamp has a period of oscillation shorter than the average recombination period of the positive and negative ions, yet not so short as to unduly reduce the effectiveness of transistor switching. It has moreover been recognized that heating at least one cathode in the lamp just prior to or while the starting voltage is being applied will facilitate the ionization process."}
{"text": "1. Field of the Invention\nThe present invention relates to a Raman amplifier that amplifies a wavelength division multiplexed (WDM) signal by utilizing a Raman effect, a WDM optical communication system, and a method for controlling Raman amplification, and in particular, relates to technology for controlling pumping light in order to realize the Raman amplification coping with variations in system operating conditions.\n2. Description of the Related Art\nA WDM optical communication method that multiplexes optical signals having different wavelengths to transmit them in a single optical fiber is an effective means for realizing a more economical optical communication system of larger capacity as well as a flexible optical network. Particularly, as technology with regard to optical amplifiers has been developed in recent years, the WDM optical communication method becomes in practical use rapidly as backbone technology supporting the Internet service.\nThe WDM optical communication system enables long-distance transmission mainly by repeating signal light while amplifying it in a repeater stage using optical amplification repeating technology. As an optical amplifying means in such a WDM optical communication system, there are used a rare-earth element doped optical fiber amplifier, a Raman amplifier, and so on.\nWith regard to the Raman amplifier mentioned above, two amplification types are known: that is, a distributed parameter type and a concentrated type. The distributed parameter type is an amplification type in that pumping light is introduced into a transmission path (for example, a silica-based optical fiber and the like) of an optical communication system, to Raman amplify in a distributing manner signal light being propagated through the transmission path, so that a part of loss that occurs when the signal light passes through the transmission path can be compensated. On the other hand, the concentrated type is an amplification type in that the pumping light is introduced in a concentrated manner into an optical fiber having a small effective cross-sectional area, for example, and high non-linearity for Raman amplification.\nThe Raman amplification of both of the distributed parameter type and the concentrated type described above shows amplification characteristics having a gain peak at a frequency lower than the frequency of the pumping light by a Raman shift amount (for example, 13.2 THz in the case of silica-based medium). Therefore, it is possible to amplify signal light of an arbitrary wavelength by preparing a pumping light source of an appropriate pumping light wavelength in consideration of the shift frequency of the Raman gain, and it is also possible to Raman amplify WDM signal light having wider bandwidth arbitrarily by combining a plurality of pumping light of different wavelengths. Further, by changing a distribution ratio of pumping light power of each of the wavelengths, a profile of the Raman gain corresponding to each pumping light wavelength may also be changed so that wavelength dependence of the Raman amplified WDM signal light can be adjusted arbitrarily.\nBy the way, in the optical network system of next generation, for example, it is anticipated that operating conditions of the system such as a wavelength band of the WDM signal light and the number of signal light, a signal light level input to an optical transmission path, a type of the optical transmission path and the like may be changed dynamically. However, in the previously proposed technology regarding the Raman amplification, there is a problem in that it is difficult to flexibly cope with the dynamic change of the system operating conditions described above.\nMore specifically, for example, if the operating conditions are changed such that so-called S-band (a wavelength band of 1480 nm–1520 nm) is added to WDM signal light using C-band (a wavelength band of 1525 nm–1565 nm) and L-band (a wavelength band of 1570 nm–1620 nm), since most of the conventional Raman amplifiers have been designed individually corresponding to each band, it is possible to cope with a change of operating wavelength in one band by controlling the power, wavelength and the like of the pumping light, but it is difficult to cope with such a dynamic change across a plurality of bands using a single type of Raman amplifier.\nFurther, when the signal light of S-band is added as described above in order to achieve the wider bandwidth of WDM signal light, there is a possibility that a wavelength of pumping light used before the change of the operating conditions may coincide with or considerably approximate to a wavelength of the added S-band signal light, resulting in a problem in degradation of transmission quality of the WDM signal light.\nStill further, it is anticipated that the dynamic change of the operating conditions described above is practically carried out in stepwise along with the lapse of required time. In such a case, it is important to cope with the change of the operating conditions without affecting the service in operation. Therefore, it is necessary to establish controlling technology for optimizing a gain and gain wavelength characteristics of the Raman amplification in response to the change with time of the operating conditions."}
{"text": "FIG. 1 shows a shoe 10 to be worn on a foot of a user with a portion of the shoe 10 cut away so that the inside of the shoe 10 is partially visible. The shoe 10 includes an upper 14 and a sole 18 coupled to the upper 14. The upper 14 covers the top and sides of the user's foot, and the sole 18 covers the bottom of the user's foot and makes contact with the ground. The sole 18 typically includes an insole 22, a midsole 26, and an outsole 30 which cushion and protect the user's foot while the user makes contact with the ground. The insole 22 contacts the user's foot, the outsole 30 contacts the ground, and the midsole 26 is arranged between the insole 22 and the outsole 30. The insole 22 generally provides a comfortable surface for contact with the user's foot and is typically essentially comprised of a thin layer of a man-made material such as, for example, ethylene vinyl acetate. The midsole 26 generally provides most of the cushioning and shock absorption for the foot of the user and is typically essentially comprised of a polymer such as, for example, polyurethane, surrounding another material such as, for example, a foam, a gel, or capsules filled with air. The outsole 30 generally provides a durable surface which can sustain repeated impact and friction with the ground and is typically essentially comprised of a rubber such as, for example, carbon rubber or blown rubber.\nThe sole 18 includes a heel end 34 arranged where a user's heel is positioned when wearing the shoe 10 and a toe end 38 arranged opposite the heel end 34 where the user's toes are positioned when wearing the shoe 10. The sole 18 also includes a medial side 42 arranged closest to the user's center of symmetry when wearing the shoe 10 and a lateral side 46 arranged opposite the medial side 42 farther from the user's center of symmetry when wearing the shoe 10.\nTurning now to FIG. 2 and FIG. 3, schematic drawings of a user's foot 50 are shown including a heel 54, toes 56, an arch 58, a medial side 60, and a lateral side 62. FIG. 2 depicts a perspective lateral side view of the bone structure of the foot 50, and FIG. 3 depicts a bottom view of the foot 50 including a plurality of regions located relative to the heel 54, toes 56, arch 58, medial side 60, and lateral side 62. A calcaneus region 66 (shown in FIG. 3) on the bottom of the foot 50 is located substantially beneath a calcaneus bone 68 (shown in FIG. 2) of the user, near the heel 54. A talus region 70 (shown in FIG. 3) on the bottom of the foot 50 is located substantially beneath a talus bone 72 (shown in FIG. 2) of the user, between the heel 54 and the arch 58. A longitudinal arch region 74 (shown in FIG. 3) on the bottom of the foot 50 is located substantially beneath a navicular bone 76, a cuboid bone 78 and cuneiform bones 80 (shown in FIG. 2) of the user, near the arch 58. A metatarsal region 82 (shown in FIG. 3) on the bottom of the foot 50 is located substantially beneath metatarsal bones 84 (shown in FIG. 2) of the user, between the arch 58 and the toes 56. A ball of the foot region 86 (shown in FIG. 3) on the bottom of the foot 50 is located substantially beneath the metatarsal-phalangeal joints 88 and sesamoids 90 (shown in FIG. 2) of the user, between the arch 58 and the toes 56 and closer to the medial side 60 than the lateral side 62. A toe region 92 (shown in FIG. 3) on the bottom of the foot 50 is located substantially beneath phalangeal bones 94 (shown in FIG. 2) of the user, near the toes 56.\nWhen propelling himself on his feet, the user applies different amounts of pressure at different times to the various bones in each foot 50 during what is known as a gait cycle. For example, during a typical walking motion, the gait cycle begins when the user first contacts the ground with the heel 54 of his foot 50, thereby applying pressure to the calcaneus bone 68. As the user shifts his weight forward on his foot 50, he applies less pressure to the calcaneus bone 68 and begins to apply pressure to the talus bone 72, the navicular bone 76, the cuboid bone 78, and the cuneiform bones 80. As the user begins to propel himself off his foot 50, he applies less pressure to the talus bone 72, the navicular bone 76, the cuboid bone 78, and the cuneiform bones 80 and begins to apply pressure to the metatarsal bones 84. As the user propels himself forward, he applies pressure along the metatarsal bones 84 and to the metatarsal-phalangeal joints 88 and sesamoids 90. Finally, as the user begins to toe off and end contact with the ground, he applies less pressure to the metatarsal-phalangeal joints 88 and sesamoids 90 and applies pressure to the phalangeal bones 94. Finally, to toe off, the user applies pressure to the phalangeal bones 94 to propel forward. The user then lifts his foot 50 into a leg swing, and places it down in a location forward relative to where he lifted it. When the user places his foot 50 down again, he first contacts the ground with the heel 54, beginning a new cycle of the walking motion.\nMany styles of forward propulsion, including many styles of walking and running, apply a gait cycle substantially similar to that described above. In some styles of forward propulsion, such as, for example, sprinting or shuffling, different amounts of pressure are applied to different portions of the foot 50 for different amounts of time. Additionally, the particular amounts of pressure applied to different portions of the foot 50 can vary from one individual to another. For example, some individuals apply more pressure to the medial side 60 than the lateral side 62 as they progress through the gait cycle. This particular application of pressure is known as pronation. In contrast, some individuals apply more pressure to the lateral side 62 than the medial side 60 as they progress through the gait cycle. This particular application of pressure is known as supination. Additionally, some individuals apply more pressure to their heels 54 when contacting the ground and some contact the ground with a portion of their feet nearer to the arch 58.\nShoes are designed to support and protect the feet of users during gait cycles to provide comfort and to promote efficient propulsion. However, due to differences between individuals in both foot anatomy and personal gait cycle style, some shoes are more comfortable and useful for some users than others. Additionally, producing a shoe configured to meet the variety of needs during all stages of the gait cycle can include producing a large number of different specialized parts which must be assembled into the shoe. Production and assembly of parts are contributing factors to the cost of the shoe. In general, a shoe having a larger number of parts is more expensive to produce than a shoe having a smaller number of parts. In view of the foregoing, it would be advantageous to provide a shoe that is comfortable and useful for a user and that is inexpensive to produce. It would also be advantageous to provide a shoe with a support arrangement that can be easily customized to meet the unique needs of various foot anatomies and individual gait styles."}
{"text": "1. Field of the Invention\nThe present invention relates to a developer for developing latent electrostatic images for use in electrophotography, electrostatic printing and the like.\n2. Discussion of Background\nTwo-component dry developers comprising carrier particles and toner particles have been conventionally known. In such a two-component dry developer, finely-divided toner particles are held on the surface of compartively large carrier particles by the electric force generated by the friction between the carrier particles and the toner particles. When the two-component dry developer is caused to come close to a latent electrostatic image, the attraction force generated between the toner particles and the latent electrostatic images overcomes the bonding force between the toner particles and the carrier particles, so that the toner particles are caused to be deposited on the latent electrostatic images. As a result, the latent electrostatic image is developed with the toner particles to a visible toner image. Therefore, the two-component dry developer is used as the toner particles are supplied thereto from time to time while in use in compensating for the toner particles used.\nAs the materials for a carrier for two-component dry developers, metallic oxides such as magnetite and ferrite are widely used. This is because such metallic oxides have a smaller apparent density than that of an iron powder carrier, so that when such metallic oxides are used as the materials for a carrier, the weight of the two-component dry developer can be reduced. Furthermore, such metallic oxides have an advantage over other materials that when such a metallic oxide is used as a carrier for a two-component dry developer, the stirring resistance of the two-component developer in a development unit is smaller than the stirring resistance of other materials employed in the carrier.\nIn addition, such metallic oxides have a smaller residual magnetic flux density and a smaller anti-magnetization force than those of an iron powder carrier, and accordingly have a smaller hysteresis loop area than that of an iron powder carrier. Furthermore, such metallic oxides have the characteristics that initial characteristics are always maintained against magnetic reversion and magnetization hysteresis.\nSince magnetite and ferrite are oxides, they are chemically stable and hardly chemically changed in contact with ozone, NO.sub.x and the like, which are formed within a copying machine.\nA carrier comprising an oxide such as ferrite or magnetite, however, has the shortcoming that a so-called spent phenomenon that a toner film is formed on the surface of carrier particles takes place by the heat generated while in use by the collision among developer particles during high speed development or during the process of making a number of copies, or by a mechanical collision between developer particles and members for a development unit while in use. Once such a spent phenomenon takes place, the charging performance of the carrier is decreased with time while in use. As a result, the toner particles are scattered and toner particles are deposited on the background of images.\nIn order to prevent the occurrence of such a spent phenomenon, methods of coating the surface of the core particles of carrier particles with a variety of resins have been proposed. However, none of them is satisfactory for use in practice. To be more specific, carrier particles coated with styrene-methacrylate resin or styrene polymer have excellent charging characteristics, but the critical surface tension thereof is relatively high, so that the above-mentioned spent phenomenon takes place during a repeated copy making operation and the life of the developer is shortened.\nExamples of a conventional negative charge controlling agent include metal complex salts of monoazo dyes, nitrohumic acid and salts thereof, sulfonated copper phthalocyanine pigments, nitro-group- or chlorine-introduced styrene oligomers, chlorinate paraffin, and melamine resin. These compounds have a complicated structure, and therefore are unstable in the properties.\nWhen such negative charge controlling agents are kneaded with application of heat thereto, they are easily decomposed and are caused to deteriorate, so that the charge controlling performance is lowered. Furthermore, the chargeability of many of such charge controlling agents is changed by the ambient conditions thereof.\nThere is a case where when a toner comprising such a conventional charge controlling agent is used for an extended period of time, the toner is deposited in the form of a film on a photoconductor because of its improper chargeability.\nFor instance, Japanese Laid-Open Patent Application 61-223753 discloses toners comprising aromatic hydroxy metallic salts such as a salicylic acid chromium complex. These toners, however, have the shortcomings that the chargeability is unstable, and the charging performance is greatly changed depending upon the ambient conditions thereof.\nJapanese Laid-Open Patent Application 3-1162 discloses a method of using a fluorinated ammonium compound or iminium compound. However, when this method is employed, the charging stability differs depending upon a carrier employed, and it is difficult to obtain a sufficient charge-stabilizing effect on a non-coated carrier for use in practice by this method.\nIn the case of a styrene-acrylic copolymer coated carrier, the chargeability is stable in a continuous mixing process, but when it is repeatedly used with a toner being replenished in a development unit, the chargeability is unstable, and the initial charge-rising performance is not satisfactory for use in practice."}
{"text": "One example of a conventional indexing method for storage, display and search of video data is described in Patent Document 1. With this conventional method, a determination is made whether telop characters are displayed for each frame of input video data, a character area is extracted from the frame whose telop characters are displayed so that a process for recognizing the telop characters is executed. Further, an index file which includes the recognized result of the telop characters and ID information at the time of displaying the telop characters is generated. The ID information is ID information at the time of inputting the frame to be recognized. Further, the Patent Document 1 discloses that the ID information is occasionally the input time information of the frame.\nPatent Document 1: Japanese Patent Application Laid-Open No. 11-167583 (JP-A 11-167583) (paragraph 0002 and FIG. 2)"}
{"text": "Unmanned aerial vehicles (UAVs), or drones, are used in a number of different applications. For example, drones may be used to deliver first aid and other supplies in emergency situations, transporting goods for commercial applications, surveying land, capturing photographs and videos, police and security monitoring, and recreational use.\nUAVs often operate in environments in which the communication link is near line-of-sight to a base station or antenna on the ground, and in which free-space propagation increases interference from neighboring transmissions. These characteristics result in challenges in maintaining good link performance between UAVs and ground base stations (e.g., mobile network cell towers) that aid in navigation, and can also cause degrade performance to ground users who share the network with the UAVs."}
{"text": "This invention relates to memory systems, memory subsystems, memory modules or a system having memory devices. More specifically, this invention is directed toward memory system architectures which may include integrated circuit devices such as one or more controllers and a plurality of memory devices.\nSome contemporary memory system architectures may demonstrate tradeoffs between cost, performance and the ability to upgrade, for example; the total memory capacity of the system. Memory capacity is commonly upgraded via memory modules or cards featuring a connector/socket interface. Often these memory modules are connected to a bus disposed on a backplane to utilize system resources efficiently. System resources include integrated circuit die area, package pins, signal line traces, connectors, backplane board area, just to name a few. In addition to upgradeability, many of these contemporary memory systems also require high throughput for bandwidth intensive applications, such as graphics.\nWith reference to FIG. 1, a representational block diagram of a conventional memory system employing memory modules is illustrated. Memory system 100 includes memory controller 110 and modules 120a–120c. Memory controller 110 is coupled to modules 120a–120c via control/address bus 130, data bus 140, and corresponding module control lines 150a–150c. Control/address bus 130 typically comprises a plurality of address lines and control signals (e.g., RAS, CAS and WE).\nThe address lines and control signals of control/address bus 130 are bussed and “shared” between each of modules 120a–120c to provide row/column addressing and read/write, precharge, refresh commands, etc., to memory devices on a selected one of modules 120a–120c. Individual module control lines 150a–150c are typically dedicated to a corresponding one of modules 120a–120c to select which of modules 120a–120c may utilize the control/address bus 130 and data bus 140 in a memory operation.\nHere and in the detailed description to follow, “bus” denotes a plurality of signal lines, each having more than two connection points for “transceiving” (i.e., transmitting or receiving). Each connection point electrically connects or couples to a transceiver (i.e., a transmitter-receiver) or one of a single transmitter or receiver circuit\nWith further reference to FIG. 1, memory system 100 may provide an upgrade path through the usage of modules 120a–120c. A socket and connector interface may be employed which allows each module to be removed and replaced by a memory module that is faster or includes a higher capacity. Memory system 100 may be configured with unpopulated sockets or less than a full capacity of modules (i.e., empty sockets/connectors) and provided for increased capacity at a later time with memory expansion modules. Since providing a separate group of signals (e.g., address lines and data lines) to each module is avoided using the bussed approach, system resources in memory system 100 are efficiently utilized.\nU.S. Pat. No. 5,513,135 discloses a contemporary dual inline memory module (DIMM) having one or more discrete buffer devices. In this patent, the discrete buffer devices are employed to buffer or register signals between memory devices disposed on the module and external bussing (such as control/address bus 130 in memory system 100). The discrete buffer devices buffer or register incoming control signals such as RAS, and CAS, etc., and address signals. Local control/address lines are disposed on the contemporary memory module to locally distribute the buffered or registered control and address signals to each memory device on the module. By way of note, the discrete buffer devices buffer a subset of all of the signals on the memory module since data path signals (e.g., data bus 140 in FIG. 1) of each memory device are connected directly to the external bus.\nIn addition to the discrete buffer device(s), a phase locked Loop (PLL) device may be disposed on the contemporary DIMM described above. The PLL device receives an external clock and generates a local phase adjusted clock for each memory device as well as the discrete buffer devices.\nModules such as the DIMM example disclosed in U.S. Pat. No. 5,513,135 feature routed connections between input/outputs (I/Os) of each memory device and connector pads disposed at the edge of the module substrate. These routed connections introduce long stub lines between the signal lines of the bus located off of the module (e.g., control address bus 130 and data bus 140), and memory device I/Os. A stub line is commonly known as a routed connection that deviates from the primary path of a signal line. Stub lines commonly introduce impedance discontinuities to the signal line. Impedance discontinuities may produce undesirable voltage reflections manifested as signal noise that may ultimately limit system operating frequency.\nExamples of contemporary memory systems employing buffered modules are illustrated in FIGS. 2A and 2B. FIG. 2A illustrates a memory system 200 based on a Rambus™ channel architecture and FIG. 2B illustrates a memory system 210 based on a Synchronous Link architecture. Both of these systems feature memory modules having buffer devices 250 disposed along multiple transmit/receive connection points of bus 260. In both of these examples, the lengths of stubs are significantly shortened in an attempt to minimize signal reflections and enable higher bandwidth characteristics. Ultimately however, memory configurations such as the ones portrayed by memory systems 100, 200 and 210 may be significantly bandwidth limited by the electrical characteristics inherent in the bussed approach as described below.\nIn the bussed approach exemplified in FIGS. 1, 2A and 2B, the signal lines of the bussed signals become loaded with a (load) capacitance associated with each bus connection point. These load capacitances are normally attributed to components of input/output (I/O) structures disposed on an integrated circuit (IC) device, such as a memory device or buffer device. For example, bond pads, electrostatic discharge devices, input buffer transistor capacitance, and output driver transistor parasitic and interconnect capacitances relative to the IC device substrate all contribute to the memory device load capacitance.\nThe load capacitances connected to multiple points along the length of the signal line may degrade signaling performance. As more load capacitances are introduced along the signal line of the bus, signal settling time correspondingly increases, reducing the bandwidth of the memory system. In addition, impedance along the signal line may become harder to control or match as more load capacitances are present along the signal line. Mismatched impedance may introduce voltage reflections that cause signal detection errors. Thus, for at least these reasons, increasing the number of loads along the bus imposes a compromise to the bandwidth of the memory system.\nIn an upgradeable memory system, such as conventional memory system 100, different memory capacity configurations become possible. Each different memory capacity configuration may present different electrical characteristics to the control/address bus 130. For example, load capacitance along each signal line of the control/address bus 130 may change with two different module capacity configurations.\nAs memory systems incorporate an increasing number of memory module configurations, the verification and validation of the number of permutations that these systems make possible may become increasingly more time consuming. Verification involves the confirmation of operation, logical and/or physical functionality of an IC by running tests on models of the memory, associated devices and/or bus prior to manufacturing the device. Validation involves testing the assembled system or components thereof (e.g., a memory module). Validation typically must account for a majority of the combinations or permutations of system conditions and possibilities which different memory configurations (e.g., 256 Mbyte, 1 Gbyte . . . ) present including signaling electrical characteristics (e.g., impedance, capacitance, and inductance variations), temperature effects, different operating frequencies, different vendor interfaces, etc, to name a few. Thus, as the number of possible memory configurations increase, the test and verification time required also increases. More time required to test a system often increases the cost of bringing the system to market or delays a product introduction beyond an acceptable window of time to achieve competitiveness.\nThere is a need for memory system architectures or interconnect topologies that provide cost effective upgrade capabilities without compromising bandwidth. Using conventional signaling schemes, the bussed approaches lend efficiency towards resource utilization of a system and permits module interfacing for upgradeability. However, the bussed approach may suffer from bandwidth limitations that stem from the electrical characteristics inherent in the bus topology. In addition, impedance along a signal line may be increasingly more difficult to control with increased connection points along a signal line, introducing impedance mismatch and signal reflections. Utilizing the bussed approach in implementing an upgradeable memory system introduces many possible electrical permutations and combinations with each unique module configuration."}
{"text": "1. Field of the Invention\nThis invention relates, generally, to boat lifts. More particularly, it relates to a four-cable boatlift having a master/slave cable system.\n2. Description of the Prior Art\nLifting a boat from the water with a piling-mounted boatlift is problematic. Typically, a winch is used to wrap cables around drums when a boat is lifted and to unwrap the cables from the drums when a boat is lowered. Many boatlifts provide one winch for each cable, there being a winch on each piling. In addition to the expense of four winches, systems of this type also have the drawback caused by different winches operating at differing speeds. Thus, if one of the four winches operates slower or faster than the others, the boat will be unevenly lifted.\nSome systems therefore use two winches, one for the inboard (dockside) end of the boatlift and one for the outboard end thereof. Still, if the two winches do not operate at the same speed, an uneven lift occurs. Uneven lifts are undesirable because a boat supported unevenly may slide from the boatlift.\nThree piling systems are introduced by the present inventor in a co-pending U.S. patent application entitled xe2x80x9cThree Piling Boat Liftxe2x80x9d filed Jan. 3, 2001, bearing Ser. No. 09/681,095. However, that system also relies upon two independent winches.\nWhat is needed, then, is a cable transmission system that synchronizes the winding and unwinding of cables around their respective drums so that boats are lifted evenly.\nHowever, in view of the prior art considered as a whole at the time the present invention was made, it was not obvious to those of ordinary skill in the pertinent art how such need could be fulfilled.\nThe long-standing but heretofore unfulfilled need for a synchronous cable transmission system for lifting boats is now met by a new, useful, and nonobvious invention.\nA three piling boatlift embodiment of the invention includes an upstanding inboard piling and a pair of upstanding outboard pilings disposed in transversely spaced apart relation to the inboard piling and in longitudinally spaced apart relation to one another. A truncate rigid plate is positioned in closely spaced relation to the inboard piling and in substantially parallel relation to the pair of outboard pilings. A first beam has an inboard end secured to a first end of the truncate rigid plate and an outboard end disposed in closely spaced relation to a first outboard piling of the pair of outboard pilings. A second beam has an inboard end secured to a second end of the truncate rigid plate and an outboard end disposed in closely spaced relation to a second outboard piling of the pair of outboard pilings.\nA cable drum assembly is mounted to the inboard piling near an upper end thereof in vertically spaced relation to the truncate rigid plate and includes a longitudinally disposed master cable drum and a slave cable drum disposed in transversely spaced apart, parallel relation to the master cable drum. The master cable drum and the slave cable drum are disposed in a common horizontal plane. A motor means rotates the master cable drum. An interconnecting means translates rotation of the master cable drum into simultaneous and corresponding rotation of the slave drum.\nA first master cable extends from a first end of the master cable drum to the first piling of the pair of outboard pilings and is adapted to lift the outboard end of the first beam when the master cable drum is rotated in a first direction and to lower the outboard end of the first beam when the master cable drum is rotated in a second direction. A second master cable extends from a second end of the master cable drum to the second piling of the pair of outboard pilings and is adapted to lift the outboard end of the second beam when the master cable drum is rotated in a first direction and to lower the outboard end of the second beam when the master cable drum is rotated in a second direction.\nA first slave cable extends from a first end of the slave cable drum to a first end of the truncate rigid plate and is adapted to lift the inboard end of the first beam when the slave cable drum is rotated in a first direction and to lower the inboard end of the first beam when the slave cable drum is rotated in a second direction. A second slave cable extends from a second end of the slave cable drum to a second end of the truncate rigid plate and is adapted to lift the inboard end of the second beam when the slave cable drum is rotated in a first direction and to lower the second beam when the slave cable is rotated in a second direction.\nThe weight of a boat lifted and lowered by the boatlift is therefore distributed substantially equally between the first and second master cables and the first and second slave cables. Moreover, all of the cables are in synchronization with one another so that the boatlift is lifted and lowered in a substantially level plane.\nThe primary object of this invention is to provide a synchronous cable transmission system where all of the cables in the system are synchronized with respect to one another to avoid uneven boat lifting.\nA more particular object is to provide a boatlift where a single winch drives a single master cable drum and all parts connected thereto.\nThese and other important objects, advantages, and features of the invention will become clear as this description proceeds.\nThe invention accordingly comprises the features of construction, combination of elements, and arrangement of parts that will be exemplified in the description set forth hereinafter and the scope of the invention will be indicated in the claims."}
{"text": "1. Technical Field\nOne or more embodiments of the present invention relate generally to the field of computer systems. More particularly, one or more embodiments of the present invention relate to the field of add-compare-select operations performed by computer systems.\n2. Description of Related Art\nTo address the needs of multi-media, communications, and graphics applications, computer systems have been designed to support one or more digital signal processing techniques to process, for example, analog data, video data, and/or audio data. Pattern recognition is one digital signal processing technique that may be used, for example, for error correction in communications applications, for image recognition applications, for speech recognition applications, and for handwriting recognition applications. Digital signal processing techniques typically use sophisticated algorithms that perform the same operations, such as an add-compare-select operation for example, on a relatively large number of data in units of bytes, words, or doublewords, for example.\nA typical computer system supports add-compare-select operations with multiple instructions and one or more arithmetic logic units. As one example, the Intel® 32-bit Architecture as defined by Intel® Corporation of Santa Clara, Calif. supports add-compare-select operations with various instructions to add packed data and with various instructions to compare packed data."}
{"text": "The present invention relates to improvements in radar systems and more particularly to improved techniques for eliminating or cancelling interference introduced into the side lobes of an ECCM radar antenna from multiple interfering signal sources.\nRadar antenna systems, including those adapted for ECCM techniques, have characteristics that include a main lobe for receiving the desired information, and a plurality of side lobes at various angles relative to the main lobe. Due to the nature of an antenna, information received in a side lobe is indistinguishable from information received in the main lobe, and thus renders the antenna highly susceptible to interference from unwanted signals or information. The problem is particularly acute in radar systems where the presence of side lobes makes it possible for a single noise jammer to be effective against a radar from any angle of azimuth. The problem becomes even more acute when multiple interference or jamming sources are used against a radar and directed from a variety of directions simultaneously. Side lobe cancellation is a fundamental approach to eliminating interference in received signals, and has been used successfully to eliminate the interference introduced from a single jamming source. One such system uses an adaptive side lobe canceller system well known in the art. In general, such systems use a signal received by an auxiliary omnidirectional antenna to cancel the interference signal received in the side lobe of the primary directional antenna. In ECCM radar systems where the desired signal is usually broadband, such as a short radar pulse, chirp pulse, or spread spectrum signal, and interference is caused by a narrow band signal, a channelized frequency selective limiter circuit is utilized for suppressing a cw or narrow band interfering signal. However, in such prior art broad band ECCM systems there is no provision for responding to wide band interfering signals which are unaffected by the frequency selective limiter circuit. Accordingly, the present invention has been developed to provide an ECCM radar system wherein both wide band and narrow band interfering signals are readily cancelled or eliminated."}
{"text": "This invention is generally concerned with compositions and methods of controlling animals and, more particularly, is directed to a novel composition for repelling animals and, in particular, dogs and cats.\nFor reasons of health and convenience, it has frequently been found to be desirable to discourage animals from frequenting certain areas. For example, garbage receptacles become both an unpleasant chore to handle and a serious potential health hazard after being ravaged by domestic animals, such as dogs and cats or by non-domestic animals, such as mice, rats, coyotes, wolves, or the like.\nChemical agents are commonly employed to discourage such animals from approaching those areas from which mankind finds it desirable to exclude them, but, while there are many chemical compounds which would effectively repel ravaging animals, there are two restrictions which severely limit the number of chemicals actually available for use. The first restriction is toxicity: the compound used must be substantially, and preferably completely, non-toxic to mammals since they will frequently be used in proximity to small children and household pets, and the probability of contact and ingestion is high. The second restriction lies with the odor of the compound: if the agent is such that it is repulsive or even unpleasant to humans, then it will not be suitable for use in populated areas. Among the relatively few compounds which meet these strictures and are known to have at least some ability to repulse animals are methyl nonyl ketone, see U.S. Pat. No. 3,474,176 and Canadian Pat. No. 978,475; and cinnamaldehyde, see Journal of Wildlife Management, 40 (1): 1976 pp 145-150 and for combinations of the two compounds as an animal repellent see U.S. Pat. No. 4,169,898.\nIt has been discovered that certain common materials, e.g. mineral oil, employed as a vehicle for olfactory animal repellents, such as cinnamic aldehyde and methyl nonyl ketone, do not dissolve the repellent. Thus, as soon as such a mixture is sprayed onto, for example, a bag, separation occurs. In addition, mineral oil is quite soluble to polyethylene, a common bag material, and hence, it is absorbed readily by the bag. As a result of this effect, there is no vehicle to hold the active ingredients after a few hours of elapsed time. Consequently, the active ingredients evaporate in a short period and are not available to repel animals."}
{"text": "Heat exchangers that utilize a single-phase working fluid to transfer heat from a heat source or to a heat sink are known as single-phase heat exchangers. Single-phase heat exchangers are used in a variety of applications ranging from radiators of conventional automobiles to more exotic water-to-ammonia heat exchangers for sustaining life in outer space, e.g., aboard a space shuttle or a space station. Single-phase heat exchangers are also used in other diverse applications, such as removing waste heat from electronic devices, e.g., microprocessors, cooling fusion reactor diverters and producing slush hydrogen.\nCompact single-phase heat exchangers are particularly desirable in applications having relatively high heat fluxes. For example, the continually increasing speeds and complexity of microprocessors cause these microprocessors to generate commensurately increasing amounts of heat. Present generation microprocessors typically have heat fluxes in the range of 5 watts/cm2 to 15 watts/cm2 The next several generations of microprocessors are predicted to have much greater heat fluxes, e.g., on the order of 50 watts/cm2 to 200 watts/cm2 or more. One type of compact heat exchanger contemplated for high flux heat transfer applications is what has become known as a normal-flow heat exchanger (NFHX). Specific embodiments of NFHXs have been previously disclosed by the present inventor, e.g., in U.S. Pat. Nos. 5,029,638 and 5,145,001. An NFHX is desirable for applications such as microprocessor cooling because it provides: (1) a single phase heat exchanger having a high surface heat flux capability; (2) a compact heat exchanger in which the working fluid experiences a generally small pressure drop as it passes through the heat exchanger; and (3) a small and lightweight heat exchanger having a high thermal transfer efficiency.\nFIG. 1 shows one embodiment of an NFHX 20 as taught in the aforementioned patents. NFHX 20 includes a heat-transfer surface 22 for thermally communicating with a heat source or sink (not shown). For example, heat-transfer surface 22 may be thermally coupled to a microprocessor for removing waste heat from the microprocessor. NFHX 20 further includes a heat-transfer element 24 comprising an inlet end 26 located opposite heat-transfer surface 22 and a plurality of closely spaced plates 28. The spaces between plates 28 define a plurality of passageways 30 that are substantially normal to heat-transfer surface 22. A plurality of outlet manifolds 32 are located adjacent, and parallel, to heat-transfer surface 22. Outlet manifolds 32 are spaced from one another and each intersects each passageway 30 to permit a working fluid 34 to flow from the passageways into the outlet manifolds.\nDuring use, working fluid 34 flows from a source (not shown) into passageways 30 via inlet end 26, and then through passageways 30 to outlet manifolds 32 in a direction substantially normal to heat-transfer surface 22, and out of heat-transfer element 24 through the outlet manifolds in a direction substantially parallel to the heat-transfer surface. As working fluid 34 flows through heat-transfer element 24, it gains, or loses, most of its heat while flowing through passageways 30. The heat is transferred to, or from, working fluid 34 via plates 28, which are in thermal communication with heat-transfer surface 22. It is the fact that the majority of heat transfer to or from working fluid 34 occurs in passageways 30 that are normal to heat-transfer surface 22 that NFHX 20 gets its name.\nUnfortunately, conventional NFHXs have a number of shortcomings. For example, for NFHX 20 to provide highly efficient heat removal, the spacing between plates 28, and, hence the depth of normal-flow passageways 30, must be very small, preferably less than 0.10 mm. Relatedly, to increase the heat transfer efficiency of NFHX 20, it is desirable to make plates 28 relatively thin, e.g., on the order of 0.25 mm or less. Accordingly, the dimensions of plates 28, passageways 30 and outlet manifolds 32 must be precisely controlled. Typically, the elements of an NFHX are fabricated using techniques such as traveling-wire EDM and traditional machining. However, these techniques are typically not capable of, or are impractical for, forming elements having the small dimensions necessary for producing NFHXs capable of handling the high heat fluxes of the next generation of microprocessors.\nIn addition, the configuration of passageway 30 and outlet manifolds 32 produces performance limitations that limit conventional NFHX 20 from achieving the needed high flux heat transfer capacity. For example, the location of outlet manifolds 32 proximate to heat-transfer surface 22 interferes with the direct conduction of heat between the heat-transfer surface and plates 28. Outlet manifolds 32 extend the length of NFHX 20 and, therefore, interrupt thermal conduction from heat-transfer surface 22 to plates 28 along the entire length of the NFHX. This interruption increases the thermal resistance between heat-transfer surface 22 and working fluid 34 in passageways 30, For this reason, it is desired to make the outlet manifolds 32 as small as possible in cross-sectional dimension, but this has the detrimental effect of increasing the pressure drop therewithin. In addition, pressure drop within outlet manifolds 32 creates non-uniform flow distribution of working fluid 34 over the cooling area of plates 28, thereby reducing the heat transfer efficiency of NFHX 20.\nThe normal-flow heat exchanger of the present invention, however, overcomes these and other shortcomings of conventional normal-flow heat exchangers."}
{"text": "1. Field of the Invention\nThe present invention is related generally to voltage regulators and more specifically to AC voltage stabilizers capable of providing a voltage regulated to within desired limits at high-power levels.\n2. Description of the Prior Art\nWith the proliferation of input voltage-sensitive loads, such as computers, there is an increasing demand for AC line voltage stabilizers. These devices are connected between the utility line and the voltage-sensitive load. They typically provide a .+-.1% output or load voltage from a .+-.17% input or utility voltage.\nFor loads up to approximately 20 KVA there is a number of methods for providing the required stabilization. The vast majority of applications can be satisfied using a constant voltage transformer or one of its derivatives. However, for powers up to the 500 KVA range the constant voltage transformer is no longer viable. Alternative methods such as motor alternator sets are expensive, heavy, and pose special siting and maintenance problems. Other methods such as motor-driven variacs are generally too slow in operation; tap changing transformers may not provide sufficiently tight control.\nSeveral static methods are available for providing the required regulation. Representative of these is U.S. Pat. No. 3,435,331. Disclosed therein is a voltage regulator utilizing a gapped booster transformer and a gapped filter transformer each having a winding comprised of a first and second portion. The first portion of each of the windings of the two transformers is connected in series with the source and the load. The first and second portions of the winding of the filter transformer are connected in series with a harmonic filter circuit across the turns of the booster transformer. The conduction of current through the second portion of the booster transformer is controlled by a pair of inverse parallel-connected silicon control rectifiers. The rectifiers are fired by a control circuit. This patent is characteristic of the prior art in that little or no protection is provided for the silicon control rectifiers by way of limiting the current flowing therethrough. It is also characteristic of the prior art in its use of harmonic filters for filtering the output wave form."}
{"text": "1. Field of the Invention\nThis invention relates generally to toys and more particularly to a means for supporting a toy craft.\n2. Description of the Prior Art\nThe prior art toy vehicles and tracks on which they operate are constructed and assembled for use on tables or the floor which interfere with other uses of the room in which these toys are used. Consequently, the use of such toys must be temporary and from time to time the tracks are dismantled and stored until there is a desire to use the toy vehicle and track again; otherwise, the room or portion of the room cannot be used for other purposes. Another disadvantage of the prior art devices to which this invention pertains is the fact that the crafts or vehicles operated on the track are not interchangeable since the vehicles are especially constructed to connect to the track by means of a retaining wheel or pin and slot arrangement which prevent the vehicles from leaving the track."}
{"text": "1. Field of the Invention\nThe invention relates to a polymeric composition having improved resistance to environmental flex cracking, solvent attack, abrasion and ozone attack. More particularly, the invention relates to a polymeric composition containing radial monoalkenyl arene conjugated diene block copolymers and an ethylene-vinyl acetate copolymer along with various other components.\n2. Description of the Prior Art\nThe use of certain block copolymers and their compounds in a number of end-uses and especially footwear, is rapidly expanding. For example, in the footwear industry, certain block copolymer compounds are being used in large volume in the manufacture of a wide variety of shoes including canvas shoes, sport shoes, and dress shoes. The block copolymers may be utilized for this purpose for soling, foxing, innersoles and the like. These compositions are disclosed in Hendricks et al, U.S. Pat. No. Re. 28,236. In general, these shoe component compositions usually comprise a mixture of block copolymer, extending oil, polystyrene and possibly a filler.\nFor the most part, these block copolymer compositions have proven to be particularly advantageous in many respects, including flexibility, processability, lack of vulcanization requirement and wet slip. However, in certain respects, shortcomings have been noted. For example, commercial compounds should be free from delamination and free from elephant hide. For some important commercial applications, it is also important that the shoe soles also possess solvent resistance and oil resistance. Improved footwear compositions have been made, such as those compositions disclosed in U.S. Pat. Nos. 4,216,131 and 4,216,132. However, the prior art compositions still do not possess sufficient solvent resistance. A novel composition with unexpectedly improved properties has now been found."}
{"text": "1. Field of the Invention\nThe present invention relates to an x-ray diagnostic installation having a portable x-ray apparatus, including an image pick-up device for generating image data, and image electronics.\n2. Description of the Prior Art\nIn the operation of a portable x-ray apparatus such as a C-bend apparatus, it is customary to accommodate the image electronics in a truck, connected to the x-ray apparatus via a fixed cable. Thus, the image electronics must also be proximally placed in relation to the x-ray apparatus, or excessive or disturbing cable lengths will be necessary to connect the x-ray means to the image electronics."}
{"text": "1. FIELD OF THE INVENTION\nThis invention relates to electrophotography and more particularly to a method for controlling the intensity of exposure of an electrophotographic continuous tone film to accurately reproduce a desired optical density range.\n2. DESCRIPTION OF THE PRIOR ART\nElectrophotographic image reproduction systems have been in existence for a number of years. In general, such systems operate as follows. An imaging element conprising a photoconductive layer that upon exposure to actinic radiation becomes conductive allowing an accumulated charge on the element surface to selectively bleed through a conductive path is first charged with a uniform charge layer by passing such element under a source of ionizing radiation, e.g., a scorotron or other such corona charging device. The charged surface is then exposed to imagewise modulated actinic radiation, rendering the photoconductor layer conductive and discharging the accumulated charge. The term \"actinic radiation\" is construed to encompass not only photochemical activity but also the photoelectric effects described herein and the like.\nIn a continuous tone system, as contemplated herein, the amount of charge left on the imaging element surface is inversely proportional to the amount of actinic radiation received by the element. In this manner a pattern of electrostatic charges is produced on the imaging element forming a latent image corresponding to the imagewise modulated actinic radiation incident on the element. The magnitude of the electrostatic charge at any one point on the imaging element is inversely proportional to the intensity of the exposing actinic radiation.\nThe latent image may now be rendered visible by development using colored particles which preferably bear a static charge and which are attracted to the charge pattern on the imaging element. Depending on the desired result, the colored particles may bear a charge of the same polarity as the charge originally placed on the imaging element or an opposite polarity. If the charge polarities are the same and an appropriate bias electrode used the colored particles are preferentially attracted to the areas from which the original charge has been bled away, producing a \"dark\" or \"colored\" area of intensity proportional to the original exposure. If the charge polarities are opposite, then the areas that received the least exposure to actinic radiation will attract the most particles. In the first instance there is an image reversal; the light tones appear dark and the dark tones appear light. In the second instance the image tones are reproduced the same as the original.\nThe colored particles may be in dry form or may be supplied in a dispersion in a carrier liquid. Generally referred to as toners, the colored particles or dispersions are well known in the art. Liquid toners tend to produce higher image resolution and are sometimes preferred for that advantage.\nFollowing toning, the image may be viewed as such, dried, fused or transferred onto a receiving element or any combination of the above, as is well known in the art.\nIn recent years the widespread use of computers and their ability to store and manipulate large amounts of data has resulted in image handling systems that employ image enhancement in applications such as radiography, printing, etc. In radiography, for instance, a radiogram may be split into a number of digitally encoded picture elements, or \"pixels\", transmitted through telephone lines, stored on a disk, retrieved at will, contrast enhanced, and displayed for diagnostic purposes. Typical display media are cathode ray tubes, silver halide film, electrostatic display, etc.\nAt present the display of high resolution diagnostic quality images is inadequate. Cathode ray tube displays have limited resolution and dynamic range. Reproduction on a silver halide film, while offering excellent resolution and dynamic range, is expensive, usually time consuming and requires darkroom facilities. Electrophotography is very promising since it reproduces high resolution images of sufficient dynamic range rapidly without the need for dark room development and complicated chemical processes. However, in order to obtain the required diagnostic quality in the finished product the exposure intensity level must be controlled to compensate for the electrostatic charge-retaining characteristic response of an electrophotographic film and for the toner electrostatic response. To complicate matters neither the response of the film nor of the toner is linear, and both tend to vary with time, usage and/or environmental conditions.\nAccordingly, in view of the foregoing, it is believed advantageous to provide a system for the accurate reproduction of the tonal range in a continuous tone image."}
{"text": "The invention relates to an encoder particularly for high speed and high resolution flash analog-to-digital converters (ADC).\nTwo parts can be distinguished in every flash ADC: an analog and digital sections. The analog section of conventional N-bit flash ADCs consists of a reference source, chain of 2.sup.N -1 resistors and the same number of comparators. The digital section consists of an encoder for sampling comparator output signals, correcting faulty codes and providing binary output code.\nThe digital output code from the analog section of a flash ADC provided by comparators is referred to as thermometer code since only one pair of adjacent bits is 0 and 1. The code must be sampled and stored as it is unstable due to changes of the ADC input signal, offset voltage of comparators, noise, etc.\nThe huge number of the comparators in high resolution flash ADCs results in a very complex gate structure of the encoders. Certain kinds of errors caused by a false response of the comparators are not possible to correct due to technological limitations. Simplified error correction schemes result in missing codes. For instance, the flash ADCs have a tendency to miss a code every few billion conversions, i.e. several times a minute when an ADC is sampling at 100MHz."}
{"text": "1. Technical Field\nThe inventive concept relates generally to nonvolatile memory devices and, more particularly, to methods of reading data from nonvolatile memory devices, the nonvolatile memory devices, and methods of operating memory systems.\n2. Description of the Related Art\nThe individual memory cells in a nonvolatile memory device, such as a flash memory device, can store data in accordance with a plurality of threshold voltage distributions, where each respective threshold voltage distribution is assigned to a corresponding logic state for stored data. The data stored by a memory cell may be read by determining whether the memory cell is turned ON/OFF when a predetermined read voltage is applied.\nDuring (and/or following) the programming of a memory cell, its intended threshold voltage distribution may be undesirably shifted or broadened due to a number of events or conditions including (e.g., charge leakage, program disturbances, word and/or bit line coupling, temperature change, voltage change, etc. Once the threshold voltage distribution of a memory cell has been thus altered, it may become difficult to accurately read the logic state of the stored data. In some instances, the logic state can become indiscernible and a read fail may occur. Once a read fail occurs, conventional nonvolatile memory devices may perform a so-called “read retry” operation.\nIn general during the read retry, a nonvolatile memory device may iteratively perform a read operation. During each successive iteration of the read operation, the level of the applied read voltage may be sequentially increased or decreased until the read operation does not result in a read fail. However, the use of repeated read operations during a read retry operation may invariably extend the time required to successfully read data from the nonvolatile memory device."}
{"text": "Electrical stimulation of neural or nervous tissue—neurostimulation—is a well established procedure for the treatment of various neurological disorders. Neurostimulation has been successfully applied to treat diseases, such as epilepsy, depression, Parkinson's disease, migraine, and for stroke rehabilitation. However, today's neurostimulation techniques have a lot of difficulties. In Vagus Nerve Stimulation (VNS) electrodes are implanted around the vagus nerve, which carries the signal rather systemically but not targeted into the brain. The VNS technique has a rather low efficacy in successfully treating epilepsy or depression. In Deep Brain Stimulation (DBS) and Responsive Neuro Stimulation (RNS), electrodes are implanted within the cranium having direct contact with the brain. The invasiveness of both techniques is a large barrier for using these techniques. In Transcranial Magnetic Stimulation (TMS), alternating magnetic fields are utilized to induce eddy currents in parts of the brain where they cause electrical stimulation. However, TMS requires rather large equipment and can not be built into a portable device that a patient can carry for the treatment of epilepsy or other chronic diseases. In Electroconvulsive Therapy (ECT), patients receive, under anesthesia, rather large electrical shocks. ECT can not be made as a portable technique to treat chronic diseases.\nIntracranial stimulation in which the stimulating electrodes are located directly at the disease site appears to be the most promising stimulation treatment for chronic illnesses, such as epilepsy, Parkinson's, depression, and others. However, intracranial stimulation techniques are very invasive.\nConcentric ring electrodes capable of delivering stimulation to defined brain areas have been disclosed in the art (e.g., WO/2006/04479), however, for diseases such as epilepsy, all current techniques can only be used if the foci of the seizure(s) are known and static, meaning the foci do not move away from the location of the electrodes or become larger beyond the region of the location of the electrodes. A partial epileptic seizure can transform into a generalized seizure when the focus is enlarged and spreads out to the size of the entire brain. Currently, there is no neurostimulation technique available that stimulates specific targeted brain tissue areas, is adaptable to moving or dynamically shape changing foci, and is less invasive."}
{"text": "1. Field of the Invention\nThe invention relates to a method of decoding (reproducing) a moving picture (image signals) which is compressed in a determined compression method and a moving picture reproducing device which performs the method. In particular, the invention relates to a method of decoding compressed moving picture at various decoding precision levels (resolutions) and reproducing the decoded moving picture and a moving picture reproducing device which performs the method.\n2. Description of the Related Art\nA conventional video signal coding technology is prescribed in a standard, such as MPEG (Moving Picture Experts Group: which is well known as an International Standard of a moving picture compressing method). Such a coding technology is used to reduce an amount of information of moving picture, and the compressed moving picture is used in communication and data accumulation.\nA user instructs a PC (personal computer) to reproduce the compressed moving picture by using a GUI (Graphical User Interface) implemented in an Operating System running on the PC. In this event, the user may change a display size of the moving picture by using a pointing device such as a mouse.\nThe change of the display size of the moving picture is done by a graphic hardware, and is not related to a decoding process of the compressed moving picture. For example, even if the display size is reduced to fifty percent in a vertical and a horizontal direction, the decoding process is performed about all signals of the compressed moving picture. However, it is to be considered that detailed information (high frequency component) included in the decoded moving picture is lost due to the reduction of the display size.\nAlso, there are various display devices to display the decoded moving picture. In general, a single display device can display the image in a plurality of resolutions. Therefore, when the image is displayed in a low resolution or in a small display size, a size of the image may be reduced so as to match with the display size.\nHowever, in the prior art, decoding process is performed about all signals of compressed moving picture, regardless of the display size. Thus, there is a problem that decoding process of high frequency components, which are not displayed after all, is performed even when the display size is smaller than an original display size of the moving picture."}
{"text": "Cancer is a major public health problem, accounting for roughly 25% of all deaths in the United States. Though many treatments have been devised for various cancers, these treatments often vary in severity of side effects. It is useful for clinicians to know how aggressive a patient's cancer is in order to determine how aggressively to treat the cancer.\nFor example, most patients with early-stage asymptomatic prostate cancer are treated with radical prostatectomy or radiotherapy and optionally adjuvant therapy (e.g., hormone or chemotherapy), all of which have severe side effects. For many of these patients, however, these treatments and their associated side effects and costs are unnecessary because the cancer in these patients is not aggressive (i.e., grows slowly and is unlikely to cause mortality or significant morbidity during the patient's lifetime). In other patients the cancer is virulent (i.e., more likely to recur) and aggressive treatment is necessary to save the patient's life.\nSome tools have been devised to help physicians in deciding which patients need aggressive treatment and which do not. In fact, several clinical parameters are currently in use for this purpose in various different cancers. In prostate cancer, for example, such clinical parameters include serum prostate-specific antigen (PSA), Gleason grade, pathologic stage, and surgical margins. In recent years clinical parameters have been made more helpful through their incorporation into continuous multivariable postoperative nomograms that calculate a patient's probability of having cancer progression/recurrence. See, e.g., Kattan et al., J. CLIN. ONCOL. (1999) 17:1499-1507; Stephenson et al., J. CLIN. ONCOL. (2005) 23:7005-7012. Despite these advances, however, many patients are given improper cancer treatments and there is still a serious need for novel and improved tools for predicting cancer recurrence."}
{"text": "The present invention relates to a bit demultiplexor circuit for demultiplexing a serial data stream, and a delay line circuit/clock multiplying circuit, in particular for use in a multiplexor and/or demultiplexor circuit.\nFurthermore the invention also relates to a bit multiplexor circuit, including\nan internal clock generator, which by means of a reference clock generates a number of accurately mutually time delayed clock signals,\naligning means arranged to align, by means of one of said clock signals, all incoming parallel data bits with respect to the same clock,\ndelay means for delaying, by means of another one of said clock signals, a number of said data bits,\nmultiplexing means for multiplexing the aligned and delayed data bits by means of said time delayed clocks.\nThrough International Patent application Ser. No. PCT/SE92/00809 a method for clock recovery of a digital data signal is known, wherein a number of mutually phase shifted auxiliary clock signals are used to generate a recovered clock signal for the data signal based upon the result of detection of a phase position error, if any, between the data signal and its recovered clock signal. If the phase position error is different from zero and the phase position of the recovered clock signal is between the phase positions of two auxiliary clock signals, these two auxiliary clock signals are mixed for generating an adjusted recovered clock signal with the same phase position as the data signal."}
{"text": "It is common for electronic circuits have a continuous energy supply, such as a steady dc voltage, during operation. However, in some cases it may be necessary to use a power source that is transitory or transient in nature, such as the energy pulse generated by a piezoelectric generator. Piezoelectric generators and other forms of energy harvesting device are finding use in a wide range of applications where it is undesirable or impossible to provide a wired power supply or where replacement of batteries is impractical. Energy harvesting is also beneficial for the environment, by reducing the need for production and disposal of batteries, which include toxic components, and by making use of renewable energy sources and “waste” energy. An example of an application using a piezoelectric generator is in powering a wireless switch. The switch may emit a radio signal when actuated, for example, to turn on or off a device configured to receive the signal. The switch may be arranged such that actuation strikes a piezoelectric crystal, generating a voltage pulse. This pulse charges a capacitor, which in turn powers a radio circuit that emits the signal, avoiding the need for a battery or other stored power source.\nWhere power is supplied by transient energy pulse, it is known to use the energy pulse to charge a capacitor, and the energy stored on the capacitor may then be used to power an electronic circuit or system. An exemplary arrangement is shown in FIG. 1, which shows a voltage source 100 that generates a transient signal Vin, that is passed to power management circuit 110. Power management circuit 110 supplies a current ILoad to power an electronic system 120. The power management circuit 110 is shown in more detail in FIG. 2, and includes a capacitor C which acts as a storage element and is charged by the energy pulse received at Vin. A diode 220 between the capacitor C and Vin prevents a current flowing from the capacitor to Vin when the voltage across C, VC, is greater than that at Vin, such as immediately following the energy pulse. The voltage VC is supplied to Vout, which connects to the circuit or system to be powered (ILoad)."}
{"text": "1. Field of the Invention\nThis disclosure relates to improved configurations for solid-state cooling, heating and power generation systems.\n2. Description of the Related Art\nThermoelectric devices (TEs) utilize the properties of certain materials to develop a temperature gradient across the material in the presence of current flow. Conventional thermoelectric devices utilize P-type and N-type semiconductors as the thermoelectric material within the device. These are physically and electrically configured in such a manner that the desired function of heating or cooling is obtained.\nThe most common configuration used in thermoelectric devices today is illustrated in FIG. 1A. Generally, P-type and N-type thermoelectric elements 102 are arrayed in a rectangular assembly 100 between two substrates 104. A current, I, passes through both element types. The elements are connected in series via copper shunts 106 saddled to the ends of the elements 102. A DC voltage 108, when applied, creates a temperature gradient across the TE elements. TEs are commonly used to cool liquids, gases and solid objects.\nSolid-state cooling, heating and power generation (SSCHP) systems have been in use since the 1960's for military and aerospace instrumentation, temperature control and power generation applications. Commercial usage has been limited because such systems have been too costly for the function performed, and have low power density so SSCHP systems are larger, more costly, less efficient and heavier than has been commercially acceptable.\nRecent material improvements offer the promise of increased efficiency and power densities up to one hundred times those of present systems. However, Thermoelectric (TE) device usage has been limited by low efficiency, low power density and high cost.\nIt is well-known from TE design guides (Melcor Corporation “Thermoelectric Handbook” 1995 pp. 16-17) that in today's TE materials, the cooling power at peak efficiency produced by a module with ZT=0.9 is about 22% of the maximum cooling power. Thus, to achieve the highest possible efficiency, several TE modules are required compared to the number required for operation at maximum cooling. As a result, the cost of TE modules for efficient operation is significantly higher and the resulting systems are substantially larger.\nIt is known from the literature (for example, see Goldsmid, H. J. “Electronic Refrigeration” 1986, p. 9) that the maximum thermal cooling power can be written as:qCOPT=IOPTαC−½I2OPTR−KΔT,  (1)where:                qCOPT is the optimum cooling thermal power;        IOPT is the optimum current;        α is the Seebeck Coefficient;        R is the system electrical resistance;        K is the system thermal conductance;        ΔT is the difference between the hot and cold side temperatures; and        TC is the cold side temperature.Further, from Goldsmid's:        \n                                          I            OPT                    =                                                    α                R                            ⁢                              1                                  (                                                                                    ZT                        AVE                                            -                      1                                                        )                                                      =                          α                              R                ⁡                                  (                                      M                    -                    1                                    )                                                                    ,                            (        2        )            where:                Z is the material thermoelectric figure of merit;        TAVE is the average of the hot and cold side temperatures; and        M=√{square root over (ZTAVE+1)}.Substitution Equation (2) into (1) yields:        \n                              q          OPT                =                              [                                                                                ZT                    C                                                        (                                          M                      -                      1                                        )                                                  ⁢                                  (                                                                                    Δ                        ⁢                                                                                                  ⁢                        T                                                                    T                        C                                                              -                                          1                                              2                        ⁢                                                  (                                                      M                            -                            1                                                    )                                                                                                      )                                            -                              Δ                ⁢                                                                  ⁢                T                                      ]                    ⁢                      K            .                                              (        3        )            \nThe term on the right side of Equation (3) in brackets is independent of the size (or dimensions) of the TE system, and so the amount of cooling qOPT is only a function of material properties and K. For the geometry of FIG. 1, K can be written as:\n                              K          =                                    λ              ⁢                                                          ⁢                              A                C                                                    L              C                                      ,                            (        4        )            where λ is the average thermal conductivity of the N & P materials; AC is the area of the elements; and L is the length of each element.\nSince α is an intrinsic material property, as long as the ratio Lc/Ac is fixed, the optimum thermal power qOPT- will be the same. For current equal to IOPT, the resistance is:\n                                          R            C                    =                                                    R                OC                            +                              R                PC                                      =                                                                                ρ                    TE                                    ⁢                                      L                    C                                                                    A                  C                                            +                              R                PC                                                    ,                            (        5        )            where ρTE is the intrinsic average resistivity of the TE elements; ROC is the TE material resistance; and RPC is parasitic resistances.\nFor the moment, assume RP is zero, then R is constant. IOPT is constant if LC/AC is fixed. Only if the ratio Lc/Ac changes, does K and hence, qCOPT and ROC and hence, IOPT changes.\nGenerally, it is advantageous to make a device smaller for the same cooling output. An important limitation in thermoelectric systems is that as, for example, the length LC is decreased for fixed AC, the ratio of the parasitic resistive losses to TE material losses, φC becomes relatively large:\n                              ϕ          C                =                                            R              PC                                      R              OC                                .                                    (        6        )            \nThis can be seen by referring to FIG. 1C, which depicts a typical TE couple. While several parasitic losses occur, one of the largest for a well-designed TE is that from shunt 106. The resistance of shunt 106 per TE element 102 is approximately,\n                                          R            PC                    ≈                                    (                                                                    B                    C                                    +                                      G                    C                                                                                        W                    C                                    ⁢                                      T                    C                                                              )                        ⁢                          P              SC                                      ,                            (        7        )            where GC is the gap between the TE elements; BC is the TE element and shunt breadth; WC is the TE element and shunt width; TC is the shunt thickness; and PSC is the shunt resistivity.\nFor the geometry of FIG. 1, the resistance for a TE element is:\n                                          R            OC                    =                                                    P                TE                            ⁢                              L                C                                                                    B                C                            ⁢                              W                C                                                    ,                            (        8        )            where Lc is the TE element length.Thus, using Equations (7) and (8) in (6):\n                              ϕ          C                ≈                                            B              C                        ⁡                          (                                                                    B                    C                                    +                                      G                    C                                                                                        T                    C                                    ⁢                                      L                    C                                                              )                                ⁢                                    (                                                P                  SC                                                  P                  TE                                            )                        .                                              (        9        )            "}
{"text": "Existing shower head are mostly switched by linkage of a decorative cover and a top cover, the top cover is disposed with a plurality of outlet holes, the decorative cover has a plurality of outlet holes as well. Chinese patent database with publishing number CN201195132Y is provided with a multi-functional back-switching type hand shower head, which is disclosed with a shower head comprising a handle, an outlet housing and a switch device, the switch device comprises a switch rotation disk and a water diversion disk, the shower head is switched by rotating the switch device to control the outlet waterways. Although it is switched from a back side, the outlet housing is still disturbed with outlet holes, thus making the switch mechanism complicated and not attractive."}
{"text": "Compositions comprising solution styrene butadiene rubber (SSBR) are used in such items as tire treads often in combination with a second rubber such as cis-butadiene rubber (BR) or natural rubber (NR). The characteristics of the compositions such as heat build up (HBU), rolling resistance, grip, and/or tear are believed to be heavily influenced by the composition of the SSBR, e.g., the styrene and vinyl content of the incorporated butadiene. Typically, an increase in styrene content, vinyl content, or both may improve properties like grip, tear, and/or braking performance but worsen properties like HBU, rolling resistance and/or abrasion resistance.\nVarious methods and compositions have been tried in attempts to alter the resulting polymer and improve characteristics such as wet grip and abrasion balance. As EP1637353 describes, mixtures of oil extended high styrene, low vinyl SSBR with cis polybutadiene, silica, processing oil and organosilicon compounds may be used in an attempt to balance these properties. However, additional compounds are still needed which have improved grip and rolling resistance behavior when combined with natural rubber and/or carbon black as filler.\nDE3108583 proposes to polymerize styrene and butadiene in such a way that a part of the monomers with a higher ratio of styrene than the final ratio of the target polymer is put into the reaction zone containing solvent. Simultaneously, a styrene/butadiene mixture containing the same or lower weight ratio of styrene than the targeted polymer is charged over a certain time to the reaction zone. The resulting polymer is coupled with a di-functional coupling agent. The coupling efficiency is about 50%. Unfortunately, the resulting polymer has a higher styrene content on the chain ends and a lower styrene content in the middle of the polymer chain. Such styrene blocks often result in deficient HBU and rolling resistance properties.\nDE2247470 proposes using a dilithio initiator and charging the monomer mixture to the reaction zone at a higher or equal rate than the polymerization growth rate. Unfortunately, such dilithio initiators are not cost-effective and the resulting polymer often has small styrene blocks on the end of the polymer chain. Such blocks may result in insufficient rolling resistance and/or HBU.\nIn DE1300239 a faster reacting monomer is charged to a reaction zone which contains an initiator, solvent, and a monomer mixture consisting of a specific ratio of both monomers which is maintained during the polymerization. Unfortunately, the resulting polymer has a styrene content which is only slightly increased during the long polymerization time.\nAccordingly, improved compositions of SSBR compositions are needed. It would be desirable if such improved compositions exhibit sufficient or improved balance of characteristics such as rolling resistance, HBU, abrasion, grip, and/or tear. It would further be desirable if such compositions could be made in an efficient, cost-effective manner.\nAdvantageously, improved compositions of SSBR and methods of making them have been discovered that meet many of the aforementioned needs and also have other desirable characteristics.\nIn one embodiment, a composition of the present invention comprises a styrene butadiene rubber. The styrene butadiene rubber is characterized before any vulcanization by one or more of the following: (a) two or more portions in one polymer chain which are incompatible with each other; or (b) two or more glass transition temperatures that vary by at least about 6° C.; or (c) two or more δ (SBR) values which differ by at least more than about 0.65 (J/cm3)0.5; or (d) two or more portions which vary in styrene content by at least about 20 weight percent based on the total weight of rubber. The styrene butadiene rubber comprises at least a portion which is compatible with a second rubber (e.g., a rubber used in a tire tread composition) and a portion which is not compatible with said same second rubber.\nIn another embodiment, the present invention relates to a method of making a styrene butadiene rubber. The method comprises first charging styrene and butadiene into a reactor in a first molar ratio of styrene:butadiene in the presence of a monolithio initiator, a randomizer, and a solvent under polymerization conditions sufficient to obtain a first portion of styrene butadiene rubber with a monomer conversion in one embodiment of higher than about 80%, in another embodiment higher than about 88%, and in another embodiment higher than about 99%. Next, styrene, butadiene, or a mixture thereof is charged into the reactor in a second molar ratio in the presence of said first portion of styrene butadiene rubber, randomizer, and solvent under polymerization conditions sufficient to obtain a second portion of styrene butadiene rubber by polymerization to a final monomer conversion of higher than about 95%, in another embodiment higher than about 98%, and in another embodiment higher than about 99% based on the total charged monomer. The first and second molar ratios are selected so as to produce styrene butadiene rubber characterized by one or more of the following: a) two or more portions in one polymer chain which are incompatible with each other; or (b) two or more glass transition temperatures that vary by at least about 6° C.; or (c) two or more δ (SBR) values which differ by at least more than about 0.65 (J/cm3)0.5; or (d) two or more portions which vary in styrene content by at least about 20 weight percent based on the total weight of the rubber. The styrene butadiene rubber comprises at least a portion which is compatible with a second rubber (e.g., a rubber used in a tire tread composition) and a portion which is not compatible with said same second rubber."}
{"text": "Metallic materials typified by iron are commonly used in a state of being coated for the purpose of protection from thermal and chemical damage factors, thereby improving durability and imparting an attractive appearance. These coated metals likely require various properties such as corrosion resistance, contamination resistance, and heat resistance. Applications utilizing heat resistance include, for example, exhaust components for vehicles and motorcycles, cooking utensils, and air conditioners and heating units. These applications generally require resistance to temperatures of about 300 to 400° C. and, in particular, exhaust components for vehicles require resistance to temperatures of 500° C. or more.\nFor these applications, since a postcoating method of coating after forming into a predetermined shape leads to high cost and is not suitable for mass production, a precoating method of working after previously coating has been reviewed. According to the precoating method, however, working is performed after forming a film, and thus extreme workability is generally required for the coating film. For example, a thin steel sheet likely requires the workability sufficient to cause neither cracking nor peeling of the film in a 2T to 4T bending test, such that a 180° bending is performed in the state where two to four steel sheets each having the same thickness as that of a test material are laid one upon another, and the adhesion between the film capable of enduring the above test and an underlaying metal.\nA silicone resin is known as a material of a film for heat-resistant precoated steel sheet. The heat resistance of the coating film made of the silicone resin remarkably varies depending on the kind and content of the organic group to be introduced into the silicone resin. In general, as the content of the organic group increases, the resultant film is more flexible and is excellent in workability and adhesion on working. On the other hand, the content of the organic substance should be decreased so as to enhance the heat resistance of the film. Thus, the above-mentioned workability and adhesion on working of the film are deteriorated.\nJapanese Patent Application, First Publication Nos. S63-172640, H02-265742 and H08-10701 describe, as coated steel sheets which satisfy both heat resistance and workability using a silicone resin, precoated steel sheets produced by using a coating composition containing, as a main component, a silicone resin having an alkyl group, an alkenyl group, and a phenyl group introduced therein. These coating films made of the silicone resin have a workability that is typically required of the precoated steel sheet because of comparatively high content of the organic group, but have a low resistance to temperatures of 400° C. or less. Japanese Patent Application, First Publication No. H08-245922 describes a coating film made of a monomethylsilicone resin having a comparatively low content of the organic group, and the resin coating film is generally insufficient in workability.\nTo solve such a problem, Japanese Patent Application, First Publication No. 2002-234109 describes a method of using two kinds of resins having different heat resistances in combination. Japanese Patent Application, First Publication No. 2002-307606 describes a resin coating film containing a methylsilicone resin and a linear methylphenylsilicone resin. To improve the workability of the above-mentioned coating film described in Japanese Patent Application, First Publication No. H08-245922, Japanese Patent Application, First Publication No. 2002-80974 describes a method of using a coating composition containing a monomethylsilicone resin as a main component, an isocyanate, and an epoxy resin. As described above, however, satisfactory heat resistance and workability imparted by introduction of an organic component conflict with each other, and there have never been obtained those which generally satisfy both properties at a high level.\nThe heat treatment temperature should preferably be set to a high temperature so as to form a film having high heat resistance, and the heat-resisting temperature of the film tends to decrease so as to decrease the heat treatment temperature. It was difficult to obtain a precoated steel sheet having high heat resistance because of this problem."}
{"text": "The field of the invention relates to memory devices and more particularly to non-volatile semiconductor memories.\nContinuing to increase rapidly is the use of computer memory, in particular non-volatile semiconductor memory, which retains its stored information even when power is removed. A wide variety of non-volatile memories exist. A typical commercial form of non-volatile memory utilizes one or more arrays of transistor cells, each cell capable of non-volatile storage of one or more bits of data.\nNon-volatile memory is unlike volatile random access memory (“RAM”), which is also solid-state memory, but does not retain its stored data after power is removed. The ability to retain data without a constant source of power makes non-volatile memory well adapted for consumer devices. Such memories are well adapted to small, portable devices because they are typically relatively small, have low power consumption, operate quickly, and are relatively immune to the operating environment.\nIn general, the small size, low power consumption, high speed and immunity to environment are derived from the structure of the memory. In this regard, such non-volatile memory devices are typically fabricated on silicon substrates. In addition, to obtain the advantages of small size, etc., and well as reduce costs, there is a continual effort to fabricate more circuitry within a given area.\nHighly effective approaches to increase density of nonvolatile memory include monolithic three dimensional memories disclosed in Johnson et al. U.S. Pat. No. 6,034,882, Johnson et al. U.S. Pat. No. 6,525,953, Knall et al. U.S. Pat. No. 6,420,215, and Vyvoda et al. U.S. Pat. No. 6,952,043, all hereby incorporated by reference in the entirety for all purposes.\nThe fabrication of these high-density, three dimensional memory arrays presents a number of challenges. For instance, misalignment of features during fabrication results in reduced yield and becomes more problematic as feature size is reduced. For example, in the event that a photomask is improperly placed, a memory element may be short circuited during subsequent fabrication operations. Thus, alternate methods of fabrication are needed that reduce the difficulties of aligning memory elements during fabrication while permitting improved density, decreased future size, and improved yield."}
{"text": "The present invention relates generally to plasma cutting systems and, more particularly, to a controller for use with such systems.\nPlasma cutting is a process in which an electric arc is used for cutting a workpiece. Plasma cutters typically include a power source, an air supply, and a torch. The torch, or plasma torch, is used to create and maintain the plasma arc that performs the cutting. A plasma cutting power source receives an input voltage from a transmission power receptacle or generator and provides output power to a pair of output terminals, one of which is connected to an electrode and the other of which is connected to the workpiece. An air supply is used with most plasma cutters to carry and propel the arc to the workpiece and help cool the torch.\nThere are multiple ways of initiating this cutting process, for example contact starting or high frequency or high voltage starting. Generally, in contact start plasma cutters, a movable or fixed electrode or consumable serves as a cathode and a fixed or movable nozzle or tip serves as an anode. In some units, the air supply is used to force a separation of the electrode and tip to create an initial or pilot arc. In others, mechanical or electromechanical means serve to separate the contacts and generate the pilot arc. In either case, once the pilot arc is established, air is forced past the pilot arc whereby it is heated and ionized to form a plasma jet that is forced out of the torch through the opening in the nozzle. The air aids in extending the arc to the workpiece forming a cutting arc and initiating the cutting process.\nOther systems utilize a high frequency starter to initiate the pilot and cutting arc, and still others systems employ high voltage to initiate the pilot and cutting arc. In any of the arrangements, the spaced relationship or the range of movement of the cathodic component and the anodic component require precise design and maintenance to generate the pilot arc and maintain the cutting arc.\nRegardless of which system is used to generate the pilot arc, some systems provide for multiple operating modes. One such mode is an expanded metal cutting mode. Expanded metal mode (EMM) allows an operator to perform many individual plasma cuts with a single trigger activation. Such a mode allows an operator to quickly and efficiently cut “broken metal”, such as grates, mesh, screen, chain, chain-link fencing, or any metallic material separated by an air gap. During expanded metal cutting mode, the plasma torch attempts to maintain an arc in the presence of an air gap in the workpiece. To accomplish this, when a controller predicts that the cutting arc is about to collapse due to an air gap in the workpiece, the pilot arc circuit is enabled such that, as the cutting arc collapses, a pilot arc is formed. The pilot arc is maintained internally within the plasma torch to sustain an arc during the gaps in the workpiece. Once the gap in the workpiece is traversed, the cutting arc is reestablished to cut the workpiece. This process is repeated until the expanded metal cutting operation is completed and the trigger is released.\nWhile an expanded metal cutting mode allows an operator to efficiently cut expanded metal, this mode increases wear experienced by the consumable components of the plasma torch. Although this wear is normal while in the expanded metal cutting mode, when an operator desires to perform regular solid metal cutting, operating the plasma torch in expanded metal mode exposes the components of the plasma torch to unnecessary wear. That is, in solid metal mode, the plasma torch only needs to generate a pilot arc upon actuation of the torch trigger. At the end of this process, the arc is allowed to fully collapse. Therefore, performing solid metal mode cutting while in expanded metal mode can result in excess pilot arc generation.\nExpanded metal cutting can also somewhat reduce the cutting power available at the end of a cutting process. This is particularly true with systems that rely on output current error detection or systems that rely on monitoring changes in output current to predict that an imminent cutting arc collapse. Specifically, since the “error” identified by these prediction systems will be high at the end of a cut, full power may not be available to ensure a clean finish on the cut.\nNevertheless, while resulting in unnecessary wear and less accurate cuts, some operators may inadvertently operate in expanded metal mode while cutting solid metal. One attempt to resolve this problem has been to provide a physical switch on the power supply requiring an operator to mechanically switch the plasma cutting system from one mode to another. However, in dynamic work environments, operators might be required to cut several different materials in a generally consecutive manner. Such an operation would require an operator to repeatedly change the operating mode of the plasma torch. Not only would such a requirement detrimentally affect process efficiency by requiring the operator to repeatedly stop a particular cutting process, travel to the power source to effectuate the change of cutting mode, and return to the cutting process, there would be no guarantee that the operator would not forget to switch the control or would simply ignore it and continue to cut solid metal in the expanded metal mode.\nIt would therefore be desirable to design a plasma cutting system that automatically controls the plasma torch cutter between various modes."}
{"text": "Polyester resins such as polybutylene terephthalate (PBT) and polyethylene terephthalate (PET) have excellent mechanical properties, heat resistance, moldability, recyclability, and chemical resistance, and can be made into molded articles, films, fibers and the like by melt processing. In particular, PBT is widely used as a material for industrial molded articles such as connectors, relays, and switches for automobiles and electrical/electronic devices. In general, PBT is produced from terephthalic acid (TPA) or an ester-forming derivative thereof, and 1,4-butanediol (BDO).\nIn JP 2014-129513 A, a small amount of a monomethoxy endcapped poly(ethylene glycol) is added to a thermoplastic resin, and they are reacted during polymerization to improve flowability.\nIn JP 63-35824 A, a monomethoxy endcapped poly(ethylene glycol) having a high molecular weight is reacted with polyethylene terephthalate (PET) during polymerization to impart an antifouling property and wash durability.\nIn JP 2001-247681 A, a PBT master batch containing polytetramethylene glycol is prepared to impart thermal shock resistance to PBT.\n“Synthesis and characterization of poly(ethylene glycol) methyl ether endcapped poly(ethylene terephthalate)” written by Timothy E. Long, published by Macromolecular Symposia, October 2003, volume 199, issue 1, pp. 163-172 discloses a PET resin obtained by adding a monomethoxy endcapped poly(ethylene glycol) during PET polymerization.\nPBT is known to have a higher melt viscosity as its molecular weight increases. Reducing the melt viscosity reduces the shear heating during melt processing, which enables reduced thermal decomposition, lower melt-processing temperatures, and the production of molded articles of complex shape.\nHowever, in the technique in JP 2014-129513 A, the amount of the compound having a (poly)oxyalkylene group is small since a decrease in mechanical properties should be suppressed so that the reduction effect of melt viscosity is small, which is problematic.\nIn the technique of JP 63-35824 A, the molecular weight of the compound having a (poly)oxyalkylene group is high so that the molecular weight decreases in a molten state, and the reduction effect of melt viscosity is small, which is problematic.\nIn the technique of JP 2001-247681 A, a compound having a (poly)oxyalkylene group is bound to the backbone rather than to a polymer end to prepare a block copolymer. Therefore, the reduction effect of melt viscosity is hardly obtained, which is problematic.\nIn the technique of “Synthesis and characterization of poly(ethylene glycol) methyl ether endcapped poly(ethylene terephthalate)” written by Timothy E. Long, published by Macromolecular Symposia, October 2003, volume 199, issue 1, pp. 163-172, the obtained PET resin is a low molecular weight polymer and has a low melting point and low mechanical properties. In addition, the PET resin disadvantageously gels through the introduction of a branched backbone.\nIt could therefore be helpful to provide a terminal modified polybutylene terephthalate resin having a very low melt viscosity, high melt stability, and a high melting point."}
{"text": "1. Field of the Invention\nThe present invention relates to polymer-micelle complexes.\n2. Description of Related Art\nCleaning product formulations, including those which contain common antimicrobial agents such as quaternary ammonium compounds and biguanides such as chlorhexidine and alexidine, rely on surfactants and mixtures of surfactants to deliver cleaning (detergency) and antimicrobial efficacy. A key aspect of these processes is the interaction of the surfactants and antimicrobial agents with the solid surfaces of the materials being cleaned, as well as the surfaces of microorganisms, together with the effects of the formulations on the air-water interface (surface tension). Reduction of the surface tension of aqueous formulations, which is directly related to the effectiveness of the wetting of solid surfaces and hence the detergency and antimicrobial processes, can be manipulated through the use of mixtures of surfactants, as is known in the art.\nAt a molecular level surfactants and surfactant mixtures in aqueous media exhibit the ability to adsorb at the air-water, solid-water, and oil-water interfaces, and this adsorption is hence responsible for a wide range of phenomena, including the solubilization of oils in the detergency process, the changes in the properties of solids and dispersions of solids, and the lowering of the surface tension of water. Adsorption of surfactants at interfaces is generally known to increase with surfactant concentration up to a total surfactant concentration known as the critical micelle concentration (CMC). At the CMC, surfactants begin to form aggregates in the bulk solution known as micelles, in equilibrium with the monomeric species of surfactants which adsorb onto the interfaces.\nThe details of the structures and sizes of the micelles, as well as the properties of the adsorbed layers of surfactants or surfactant mixtures, depend on the details of the molecular shape and charges, if any, on the hydrophilic “headgroups” of the surfactants. Strongly charged headgroups of surfactants tend to repel each other at interfaces, opposing the efficient packing of the surfactants at the interface, and also favoring micelle structures that are relatively small and spherical. The charged headgroups of many surfactants, such as the quaternary ammonium compounds, will also introduce a counterion of opposite charge, for example a chloride or bromide ion, into formulations.\nIt is known that the nature of the counterion can affect the repulsion between charged surfactants in micelles and adsorbed layers through a partial screening of the headgroup charges from one another in surfactant aggregates like micelles. It is also well known that addition of simple electrolytes, such as sodium chloride, into aqueous solutions can also be used to increase the screening of like headgroup charges from each other, and thus is a common parameter used to adjust the properties of surfactant micelles, such as size and shape, and to adjust the adsorption of surfactants onto surfaces.\nAddition of significant amounts of simple electrolytes into many formulations, such as hard surface spray cleaners or nonwoven wipes loaded with a cleaning lotion, is undesirable due to residues left behind upon drying of the formulations. An alternative method to adjusting the properties of such formulations, including the wetting of solid surfaces and the antimicrobial activity, is to include significant amounts of volatile organic solvents such as lower alcohols or glycol ethers. Volatile organic solvents, however, are coming under increasing regulation due to their potential health effects, and are not preferred by the significant fraction of consumers who desire efficacious cleaning and disinfecting products with a minimum of chemical actives, including volatiles. In the healthcare industry, efficacious formulations comprising quaternary ammonium compounds and lower alcohols are known, but are viewed as having shortcomings in terms of the potential for irritation of confined patients. Such products pose similar risks to cleaning and clinical personnel who may be exposed to such products on a daily basis.\nThere is an increasing interest from consumers, and a known need in the healthcare and housekeeping industries, to reduce the number of microorganisms on fabrics while using familiar equipment such as washing machines. Concentrated products are required for such an application, due to the high dilution level of the product in the rinse water, typically by a factor of about 600 times dilution. In the case of formulations comprising quaternary ammonium compounds, high concentrations of the quaternary ammonium compounds in the concentrate are needed in order to ensure an adequate amount of adsorption occurs in a kinetically relevant time onto the microbes under dilution use conditions. As detailed above, it is desirable, yet very difficult, to manipulate (i.e., reduce) the CMC of the quaternary ammonium compound in such an application. Thus very high concentrations of quaternary ammonium compounds, which tend to be hazardous to the skin and eyes, are used in the concentrates, in combination with high temperatures and long exposure times.\nThus, there is an ongoing need for methods and compositions offering fine control of the properties of surfactant aggregates comprising cationic species, especially antimicrobial species such as quaternary ammonium compounds and biguanides."}
{"text": "Technical Field\nThe disclosure relates to a phosphor, and in particular it relates to an illumination device utilizing the same.\nDescription of the Related Art\nThere are many manufacturing methods for making white light-emitting diodes (WLEDs), such as (1) applying a yellow phosphor onto a blue light-emitting diode (LED) chip, (2) applying a red phosphor and a green phosphor powder to a blue LED chip, (3) mixing a red LED chip, a green LED chip, and a blue LED chip, and (4) applying a blue phosphor, a green phosphor, and a red phosphor (or different phosphor powders of different colors) onto an ultra-violet (UV) LED chip.\nThe WLED has several advantages over traditional incandescent light bulbs, such as long lifespan, low power consumption, small volume, fast response time, and good shake-resistance. As a result, it has gradually replaced traditional lighting products. However, current WLEDs still need to overcome problems in development, such as heat dissipation, insufficient brightness, and their relatively high price. In the lighting market, WLEDs are mainly used in auxiliary lighting such as flash lights, car interior lights, or decorative architectural lighting products. WLEDs are expected to replace traditional lighting products in the future to enter the mainstream of the global lighting market.\nIn addition to packaging techniques, the phosphor that is chosen is an important factor in the luminous efficiency of a light source. One of the directions of the research being conducted by photo electronic semiconductor companies and lighting companies is modifying phosphor compositions to increase phosphor conversion efficiency. The color render index of the white light generated by a yellow phosphor powder excited by a traditional single blue chip is not good, as the color saturation of an object illuminated by the white light is poor, thereby lowering its value in the commercial lighting market. Compared with this method, the application of a green phosphor and a red phosphor onto a blue LED may produce a white light with a better color render index. Accordingly, different phosphor compositions for use in a white light illumination device are called for."}
{"text": "1. Field of the Invention\nThe present invention concerns a holding device for blood collection devices.\n2. Description of the Related Art\nThe already known blood collection devices are either syringes or those that demand the use of tubes with a vacuum treatment.\nThe blood colletion devices that demand the use of tubes with a vacuum treatment reveal two problems that the invention disclosed by this application solves through the substitution of tubes with a vacuum treatment by tubes without a vacuum treatment.\nThe first problem that the tubes with a vacuum treatment reveal is that they do not allow the person who manipulates the devices to correctly control the collection pressure of the blood that is carried by the vein because said collection pressure depends on the tube volume. Said unproper control carries as a consequence other problems. These problems are the following: a) It increases the hemolisis risk and, as a consequence, it reduces the quality of the blood samples; b) It makes difficult the obtainment of blood samples whenever the veins are very thin (e.g.: patients suffering from vein difficulties and children), increasing thrombosis or vein collapse risks.\nThe second problem that the tubes with a vacuum treatment reveal is that it is easy to loose the vacuum while manipulating the devices preventing the blood from being drawn into the tube.\nThose problems appear while using, as for instance, devices disclosed by patents U.S. Pat. No. 6,024,710, EP 867.378 A2, and JP 1.291.830 A2.\nOn the other hand, the use of syringes for the collection of blood samples has two other characteristics that carry some problems that the invention that is claimed solves.\nFirstly, the use of the syringes does not reach the higher level of the necessary asepsis for the correct manipulation of the blood samples. Said devices are opened so the blood can leak during their manipulation.\nSecondly, while using the syringes, the blood is drawn into the tube only once it has been collected by using the syringes. That is to say, there are two consecutive actions, and this fact turns the quality of the blood sample risky.\nThe Applicant does not know of any other device comprising tubes without a vacuum treatment, enabling exactly the same manipulation quality offered as those devices comprising tubes with a vacuum treatment and that, in addition, solves the problems that said devices reveal and fulfills the asepsis and blood sample quality levels as the one disclosed in this application does, and has been disclosed before the priority date that is declared in this application."}
{"text": "In passive monitoring, performance measurements are obtained by carrying out passive measurements within the network. In contrast, subscriber terminal-based monitoring characterises subscriber satisfaction by carrying out performance measurements directly in the terminal of a subscriber.\nBy way of example, for passive measurement in the network, the network traffic may be captured at certain interfaces in the network, and performance indicators may be obtained by processing this information. Performance indicators are possible to define such that end-user perceived quality is well-approximated by the indicators.\nEven though subscriber-perceived performance can generally be best observed in the subscriber terminal, a passive monitoring system for carrying out measurement in the network may have the following advantages:                no specific terminals are needed;        all terminals in a live network can be observed;        cost-efficient large-scale monitoring is possible, because a few measurement points can cover a large part of a network;        not only user-perceived performance can be observed, but also the traffic composition and volume (e.g., the most popular applications, the traffic demands of typical or power users, etc.); and        by appropriate placing of the measurement points, it is possible to localize faults or find the cause of performance degradations (e.g., bottlenecks).        \nPassive monitoring-based performance/fault analysis and traffic modeling is therefore becoming an increasingly important part of network audit and management services.\nHowever, as access rates increase and prices of services fall, an increased number of subscribers start to use more bandwidth-hungry applications. Consequently, passive monitoring becomes more difficult.\nIt is an object of the present application to provide improved systems and methods for monitoring transport of data in a generalised communications system wherein data packets conform to a layered transmission protocol."}
{"text": "The present invention relates generally to the enhancing retention of substances in biological target site, and more particularly to applying a chemical barrier for prolonging agent residence time in a biological tissue.\nThe medical and veterinary fields have long used agents to treat, diagnose and prevent countless medical conditions. The substances are used to alter the body\"\"s function at the molecular level. The types of agents that are administered to biological systems include natural and synthetic drugs, biological agents, therapeutic radiation, and the like.\nThere are many complications that accompany the use of agents. The efficacy of chemical treatment depends on several variables including (1) the mechanism of chemical reaction, (2) penetration of the chemical into a target site, (3) amount of chemical that is delivered, (4) proficiency of delivery, and (5) residence time of the chemical at the target site.\nThe composition of the chemical is one factor that impacts these variables. Enormous efforts are placed on devising new and improved compounds. In particular, chemicals are formulated for providing optimal reaction mechanisms and penetration.\nThe agent\"\"s route of administration and travel to the treatment site further influence the variables to efficacy. Traditional routes of administration include intravenous, oral, inhalational, topical, transdermal, subcutaneous, intramuscular, buccal, intra-arterial, intrathecal and rectal.\nMethods of directly administering an agent are of special interest. Such local applications reduce the distance through which the agent must travel between the site of administration and the intended target place. As a result the chemical has a higher bioavailability at the target site. Unlike systemic delivery, local delivery of chemical can be made at high concentrations without causing systemic toxicity issues. Where the agent must extensively travel through the body to reach the target tissue, the agent may be metabolized in the gut, portal blood or liver prior to entry into other systemic circulation. Moreover, an agent flowing through the circulation to reach its proposed destination may be retained by blood plasma proteins instead of binding to its intended tissue protein. Direct routes are also useful for reducing potential undesirable effects, producing high local concentrations of agent at the treatment site, and continuously attending to chronic conditions.\nIn order to enhance chemical performance, medical devices are constructed and medical procedures devised to facilitate contacting the agent with its target site. For example, many devices and methods have been designed to increase agent penetration, the amount of agent delivered and the efficiency of chemical function.\nSpecial routes of administration are available for dispensing agents directly to a target tissue. The agents may be attached to a carrier that is implanted near the target site, e.g. drug delivery stents. Other local agent delivery devices include microporous balloon catheters, ultrasound-enhanced delivery through microporous balloons, intra-vascular injection catheters, iontophoresis, etc. Intra-vascular, intraoperative or intrathecal catheters may also be surgically inserted near the target organ for delivery of agents. In an exemplary case, gene therapy agents are presented to the heart through a catheter that is percutaneously introduced, such as through the femoral artery, and guided upstream into the left ventricle and epicardium for agent introduction into the pericardial space. In such procedures, access to the pericardium may also be gained intra-vascular or through a thoracotomy.\nOnce an agent is administered, the chemicals progress throughout the body to and from the target site via the bloodstream, diffusion and active transport systems. Most fluids move freely through the circulatory system and pass to various areas of the body through capillaries. There are some natural barriers that restrict movement of molecules throughout the body, such as the blood-brain barrier for hindering the passage of selective chemicals to the brain. Most capillaries, however, have vessel walls with cells arranged in a manner that allow molecules as large as 20,000-30,000 Daltons to pass under pressure through the vessel wall.\nAfter an agent reaches its intended treatment site, it is typically absorbed by blood vessels and washed away with circulating blood flow. Agents are then usually excreted from the body by transfer via the circulation to the kidneys, liver, gastrointestinal tract, lungs, salivary glands and sweat glands.\nHowever, frequently an agent is only effective while it is in contact with the target site, e.g. bound to a receptor. Mechanisms that resist removal of an agent by circulating blood flow permit the agent to create a more dramatic outcome at the treatment site. Prolonged agent retention times at the target location enhance many treatments, especially those requiring long-term administration of agent. Moreover, since less agent is quickly washed away with these methods of increasing agent retention, a smaller amount of agent is required. Hence, the cost of the agent is reduced. Furthermore, the lower dosage of agent may decrease side effects that sometimes occur with larger amounts.\nMany medical conditions benefit from prolonged agent residence time at the target site, such as cardiac conditions. Coronary heart disease is the most common disease and cause of death in civilized, western countries. In the treatment of atheroslerotic coronary disease, percutaneous transluminal coronary angioplasty (PTCA) may be performed to reduce obstructions in a vessel. PTCA involves the introduction of a catheter with a small dilating balloon at its tip. The catheter is snaked through the arteries, often in the arm or leg, to a site where the vessel has narrowed. The balloon is inflated to increase the cross-sectional area of the vessel.\nAlthough PTCA treatment has a high success rate in widening the vessel, often the vessel re-narrows or re-closes post dilation, referred to as xe2x80x9crestenosis.xe2x80x9d In restenosis, proliferation of smooth muscle cells does not initiate until a long time (48-72 hours) after the PTCA procedure is performed. Treatment of restenosis includes delivery of anti-proliferation drugs to the site of the diseased vessel wall following PTCA and during the same surgery. Yet, the drug treatment is only mildly effective because the bloodstream removes the drug from the tissue quicker than the time it takes for the onset of restenosis. The drug is encouraged to diffuse out of the area of higher agent concentration in the vessel walls to the site of lower concentration in the internal lumen of the vessel. The flowing blood maintains this lower agent concentration inside of the vessel. In this manner, the agent is readily washed away from the target site and agent has a limited residence time in the tissue.\nCurrent approaches to increase chemical retention times involve manipulating the agent prior to its introduction into a biological system. Typically, functional molecules are pre-attached to the agent to alter the physical structure of the agent in a manner that hinders absorption of the agent after it is delivered into the blood stream. In one such technique, a polyanionic sulfate group is bound to a drug in order to impart negative charges. The negatively charged drug repels from the negatively charged walls of a capillary and resists absorption onto the vessel wall. Thus, the agent is retained within the peripheral circulation.\nSome retention systems attempt to provide for gradual release of an agent from a delivery source. Prior to administering, an agent may be coupled to a carrier that slows the release of the agent at the target site. One exemplary system binds a drug to an ionic carrier for administering to the eye. Tear fluid gradually dissolves the carrier to release the drug. Other systems include microcapsules that encapsulate the drug. The shell is biodegradable, such as by hydrolysis, to slowly release the drug.\nIn other methods, a drug is pre-incorporated into a solid polymeric material that is inserted into a vessel. Once the material reaches the target site, the material is reconfigured to form a support structure on the vessel surface. For example, the solid material may be heated to melt the material, molded to the shape of the internal vessel surface and then re-solidified. A problem with this system is that the manipulation of the material once it is in the vessel, such as by heating, may disrupt the agent and vessel.\nThere are many problems with these current systems that require the agent be structurally manipulated prior to entry into the body. A significant drawback to the use of pre-altered agents is that modification of each chemical agent is complicated and expensive. Retention mechanisms and attachment structures must be tailored for each agent of interest. The resulting change in molecular structure of the agent composition may compromise chemical action. Furthermore, pre-attachment of a delivery component to the agent results in an increased size of the agent composition. This altering of the agent may hinder penetration or delivery of the chemical.\nMoreover, carrier systems do not lengthen the time that agent molecules contact a target site, but rather provide a constant supply of agent to a site. Thus, these systems require large amounts of agent in order to lengthen the effect. These carrier systems are also limited by the amount of agent that may be administered in the vehicle.\nIn view of these limitations with current systems, there is a need for a platform for prolonging chemical residence time that may be applied universally across a variety of agents. Furthermore, a system is desired which lengthens agent effect by extending the time in which agent molecules contact a site, rather than increases the agent dosage. The mechanism should not hinder chemical reaction, penetration or delivery and cause minimal damage to the tissue. In particular, an approach for hindering reabsorption of an agent by the bloodstream without manipulating the agent prior to administering would be advantageous.\nAn agent retention system and process is provided for prolonging the residence time of an agent in, on or near a tissue. In general, the system increases agent effectiveness at a treatment site by extending the time in which agent molecules contact the site, rather than increasing the agent dosage. The tissue that benefits from the retention system may be any biological tissue having at least one surface. For example, the tissue may be a vessel, organ, tumor cells, or the like.\nA biodegradable barrier that is separate from the agent is utilized in the present invention the barrier has a binding member to couple the barrier to a tissue surface. The barrier may be a biodegradable molecule that is biocompatible with the target tissue. For example, the barrier may be a poly(anino acid), such as poly-L-lysine. The binding member permits attachment of the barrier to the tissue surface by a variety of binding mechanisms, such as ionic binding, covalent binding, electrostatic interactions, hydrogen bonding, and the like. For example, the barrier\"\"s binding member may be a cationic ion for binding to an anionic member of the tissue surface. In one embodiment of the barrier, a charged binding member is included for adhering to an oppositely charged member on an interior vessel surface and a passive platelet inhibitor.\nIn one embodiment of the agent retention method according to the present invention, a pre-treatment agent is contacted with the tissue, followed by contact with the barrier to the tissue surface. A sufficient amount of the barrier is allowed to adhere to the tissue surface to hinder diffusion of the agent through the tissue surface as compared to diffusion of the agent without the presence of the barrier.\nIn some cases, the barrier is contacted with a tissue surface prior to introduction of a post-treatment agent. In this embodiment, the post-treatment agent is contacted with the barrier and the agent attaches to the barrier. The agent is retained at the tissue for a prolonged period of time as compared to the retention of the agent without the presence of the barrier.\nIn still other configurations, the retention system is applied to a tissue having an inner tissue wall and tissue surface. The agent and biodegradable barrier are administered to a patient in a time-varied manner. A sufficient amount of the barrier\"\"s binding member is allowed to couple the barrier to the tissue surface to hinder diffusion of the agent through the tissue surface.\nThe barrier composition may be formulated by admixing a barrier having a binding member and a delivery carrier. The barrier component is present in an amount sufficient to couple to the tissue surface and to permit transport of the agent from the tissue surface at a lower rate than agent transport in the absence of the barrier composition.\nThe retention system may comprise the barrier and a conduit, e.g. a catheter, having at least one opening for administering the agent and barrier. The barrier, agent and optionally the conduit may also be included in a system for increasing agent residence time in a tissue.\nThe benefits of the agent retention system are direct in that the system may be universally applied across a variety of agents. With the present retention system, any pre-treatment agents may be retained for longer in the tissue. Agent need not be reformulated with the barrier or even compatible with the barrier. In one embodiment, the approach hinders reabsorption of an agent by the bloodstream without needing to manipulate the agent prior to administration into a patient.\nThe mechanism for agent retention according to the present invention does not hinder chemical reaction, penetration or delivery of the agent. The process further may result in only minimal damage to the tissue.\nOther features and advantages of these and other embodiments are discussed in detail below."}
{"text": "The present invention relates to a fuel supply cut control device for use in an internal combustion engine.\nIn order to improve a specific fuel consumption, a fuel supply cut control device is known in which the supply of fuel is stopped when the engine is decelerated. In such a device, the supply of fuel is stopped as soon as the decelerating operation is started, and the supply of fuel is started again when the engine speed is reduced below a predetermined fixed speed (hereinafter referred to as a resume speed). After this, when the engine is decelerated again after the engine speed is increased beyond a predetermined fixed speed (hereinafter referred to as a cut speed), the supply of fuel is stopped again until the engine speed is reduced below the resume speed. In such a fuel supply cut control device, in order to further improve a specific fuel consumption, it is necessary to reduce the resume speed and the cut speed as much as possible. However, if the resume speed and the cut speed are reduced, when the engine speed is reduced below the resume speed and then the supply of fuel is stopped, or when the decelerating operation of the engine is started under an operating state wherein the engine speed is slightly higher than the cut speed, and then the supply of fuel is stopped, a large torque fluctuation which provides a shock for a vehicle occurs. As a result, the drivability of a vehicle deteriorates. Therefore, in a conventional fuel supply cut control device, it is difficult to reduce the resume speed and the cut speed."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor apparatus and, particularly, to a semiconductor apparatus including a voltage divider that generates a divided voltage by dividing a voltage difference between a first power supply and a second power supply.\n2. Description of Related Art\nAs a recent trend, a semiconductor apparatus generates a voltage different from a power supply voltage inside the apparatus and makes an internal circuit operate with the generated voltage in order to improve the performance. An example of such a voltage is a bias voltage (or a reference voltage). The bias voltage is a voltage that serves as a reference of a circuit operation and it is required to maintain a stable voltage without depending on variations of a power supply voltage. An example of a bias voltage generator is disclosed in Japanese Unexamined Patent Application Publication No. 2003-173213.\nFIG. 15 shows a block diagram of a bias voltage generator 100 disclosed in Japanese Unexamined Patent Application Publication No. 2003-173213. Referring to FIG. 15, the bias voltage generator 100 includes a resistor element R40, an output transistor NT40, and a current control circuit 140. The resistor element R40 and the output transistor NT40 are connected in series between a power supply VCC and a power supply VSS. In the bias voltage generator 100, a node between the resistor element R40 and the output transistor NT40 serves as an output node N10. A gate and a drain of the output transistor NT40 are connected to each other. The current control circuit 140 includes a resistor element R41 and transistors NT41 and NT42. The current control circuit 140 is placed in parallel with the output transistor NT40. The current control circuit 140 increases a current 112 that flows from the output node N10 to the power supply VSS as the voltage of the power supply VCC increases. The bias voltage generator 100 can thereby allow the value of a bias voltage VBS that is output from the output node N10 to remain substantially constant or decrease with respect to an increase in the voltage of the power supply VCC.\nHowever, in the bias voltage generator 100, a voltage level that serves as a reference of the bias voltage VBS is determined by a threshold voltage (or a diode voltage) of the output transistor NT40 with respect to the power supply VSS. Some semiconductor apparatus generates a negative voltage inside the apparatus and uses it for a circuit in the apparatus in order to improve the performance. In the case of using the bias voltage generator 100 in such a case, the bias voltage VBS varies with variations of the negative voltage, and therefore the bias voltage generator 100 cannot be used for a circuit in which both the power supply VCC and the power supply VSS vary.\nDRAM (Dynamic Random Access Memory) is one of such semiconductor apparatus using a negative voltage. In DRAM, a negative voltage is generated for a substrate voltage VBB and used as a substrate voltage of a transistor of a memory cell. By using the negative voltage as the substrate voltage VBB, it is possible in the DRAM to reduce a leakage current of the memory cell and improve data storage capability (cf. “Very LSI Memory” written by Kiyoo Itoh and published from Baihukan). The negative voltage is generated typically by a charge pump. An example of a negative voltage generator that generates a negative voltage with use of a charge pump is disclosed in Japanese Unexamined Patent Application Publication No. 11-150230.\nFIG. 16 shows a block diagram of a negative voltage generator 200 disclosed in Japanese Unexamined Patent Application Publication No. 11-150230. Referring to FIG. 16, the negative voltage generator 200 includes a reference voltage generator 210, a voltage divider 220, a comparator 230, an oscillator 240, a clock buffer 250, and a negative voltage charge pump 260. The reference voltage generator 210 generates a constant voltage having a voltage value between a power supply voltage and a ground voltage and outputs it as a reference voltage 201. The voltage divider 220 divides a voltage between the reference voltage 201 and a negative voltage 204 and outputs it as a divided voltage 202. In the voltage divider 220, a voltage dividing ratio is set in such a way that the divided voltage 202 becomes a ground voltage when the negative voltage 204 reaches a voltage value to be set. The comparator 230 compares a voltage value of the divided voltage 202 with the ground voltage, and if the voltage value of the divided voltage 202 is higher than the ground voltage, it activates an oscillator control signal 203, and if the voltage value of the divided voltage 202 is lower than the ground voltage, it inactivates the oscillator control signal 203. The oscillator 240 outputs a first oscillator output signal 205 and a second oscillator output signal 206 having opposite phases when the oscillator control signal 203 is activated. The clock buffer 250 outputs a first complementary pulse signal 207 corresponding to the first oscillator output signal 205, and outputs a second complementary pulse signal 208 corresponding to the second oscillator output signal 206. The negative voltage charge pump 260 outputs the negative voltage 204 based on the first complementary pulse signal 207 and the second complementary pulse signal 208.\nThe negative voltage generator 200 has a circuit configuration in which, when the negative voltage 204 reaches a preset voltage, the divided voltage 202 becomes the ground voltage that is input as a comparative voltage of the comparator 230. In the negative voltage generator 200, the reference voltage generator 210 generates the reference voltage 201 that is not dependent on variations of the power supply voltage VCC in order to stably obtain the negative voltage 204 even when the power supply voltage VCC varies. Further, the voltage divider 220 determines a voltage dividing ratio between the reference voltage 201 and the negative voltage 204 by a resistor string. Specifically, in the negative voltage generator 200, the divided voltage 202 that reflects the voltage value of the negative voltage 204 is generated based on the stable reference voltage 201 and a fixed voltage dividing ratio. The negative voltage generator 200 can thereby stabilize the voltage value of the negative voltage 204 regardless of the power supply voltage VCC."}
{"text": "There are numerous manufactured products which are of an extended length, e.g., over eight feet in length, which are not self-supporting and which require packaging for storage and shipment. For example, thin, light-weight aluminum and plastic extrusions of eight feet or greater length are not self-supporting; and, thus when the product is lifted, it bends or kinks at one or more points along its length. With aluminum and plastic extrusions, the bending and kinking can produce a permanent defect in the product making it unusable for sale for its intended purpose. One example of such a product are thin-slat venetian blinds, commonly called mini-blinds, which are often manufactured in lengths over eight feet to span a window over eight feet in width or to be oriented vertically in a window eight or more feet high. These blinds are typically shipped in cardboard packaging which, of course, have a length equal to or greater than the length of the blind. In the process of transporting the packaging, if it is lifted at one end, the weight of the blind causes kinking of the packaging along its length and in turn kinking of the blinds. That is, the box does not have sufficient strength to support the product contained therein and in turn will deform and kink, and kink the product. Typically, for example, in a four-sided box, the sides have a tendency to deform into a parallelogram resulting in bending of the box at some point along its length."}
{"text": "More generally, the invention relates to the field of electric motors having permanent magnets such as brushless, electronically commutated DC motors and other permanent magnet motors, and in particular those configured as inner rotor motors. In general, inner rotor motors consist of a rotor arrangement which is mounted onto the motor shaft and includes one or more permanent magnets, as well as a stator arrangement, such as a stator core, which is built up of metal laminations that carry windings. The rotor arrangement is coaxially inserted into the stator arrangement. For outer rotor motors, the rotor arrangement encloses the stator.\nFIG. 8 shows the basic construction of an electric motor having a housing 114 in which a stator arrangement 118, a rotor arrangement 116 and bearings 126, 128 are accommodated to rotatably support the rotor arrangement. The stator arrangement 118 includes stacked metal laminations 155 and windings 160 and encloses an inner space into which the rotor arrangement 116 can be inserted. The rotor arrangement 116 includes the shaft 110, a back iron yoke 112 and permanent magnets 122. The bearings 126, 128 supporting the rotor arrangement can be integrated into a flange 124 in the motor housing 114.\nFIG. 8 serves to explain the basic construction of an electric motor. As explained in the opening paragraph, the invention relates to a rotor for such an electric motor, the rotor having an essentially cylindrical rotor core with a central aperture and the permanent magnets being embedded in the rotor core.\nAccording to the prior art, rotors with embedded magnets are generally known. A rotor configuration having a multi-polar design resembling a spoked wheel with radially extending embedded magnets is revealed, for example, in “Design of Brushless Permanent-Magnet Motors”, J. R. Hendershot Jr. and T J E Miller, Magna Physics Publishing and Clarendon Press, Oxford, 1994. As shown in this publication, it is known to manufacture a rotor with embedded radially extending magnets that are protected by means of a ring or a tube surrounding the rotor. The rotor in which the magnets are embedded is used as a back yoke.\nA conventional form of rotors with embedded magnets is also revealed in EP 0 691 727 A1. This publication shows a number of permanent magnets which are inserted into slots formed in the rotor allowing the permanent magnets to be inserted into the rotor from the outside. At their radially inner ends, the permanent magnets are enclosed by the material of the rotor core.\nRotors with embedded permanent magnets have the basic advantage that the magnets can be fully encapsulated so that the rotor can also come into contact with aggressive media without the magnet material needing special surface protection to prevent corrosion etc. However, the described rotor design has the disadvantage that stray flux is generated by the rotor core in the vicinity of the shaft.\nTo prevent such stray flux from arising, it has been suggested in the prior art to place a sleeve made of magnetically non-conductive or low-conductive material onto the shaft onto which the flux guide elements of the rotor core are then fixed, between which the permanent magnets in turn are embedded. Such a design is revealed, for example, in EP 0 641 059 A1; EP 0 803 962 A1; and DE 101 00 718 A1. Although this construction represents a great improvement on the prior art as described above in terms of the magnetic circuit and the distribution of magnetic flux density in the rotor, it is costly to manufacture and, due to the many individual parts, problems in the mechanical construction, such as an addition of tolerances, could arise.\nEP 0803 962 A1 additionally shows that the slots to accommodate the permanent magnets have a bridge on their outer periphery to fully protect the permanent magnets from the outside.\nWO 00/57537 describes a multi-polar permanent magnet rotor for an electric motor having embedded magnets which are disposed in such a way that a concentration of flux is produced. The permanent magnets are formed as flat cubes which are disposed like spokes radially to the rotor axis in recesses that are arranged between the flux guide elements which are fixed to the rotor. In assembling the magnets and the flux guide elements, the permanent magnets are formed as adjacent half-elements representing one pole respectively, and both the permanent magnets and the flux guide elements are attached to the shaft via a sleeve.\nAnother method of constructing a rotor having embedded magnets is shown in EP 0 872 944 A1. The magnets are arranged in a radial direction, or parallel to a radial direction, to the rotor. In EP 0 872 944, the permanent magnets are disposed in a so-called double-spoke configuration. Each of these “double magnets” consists of a pair of permanent magnets whose direction of magnetization is substantially the same. They can be arranged parallel to each other as in the cited publication or inclined at an angle to each other. This arrangement goes to improve the running performance of the electric motor and, in particular, to reduce cogging torque and torque ripple.\nOther published patents in respect of rotors with embedded magnets include GB 1,177,247; EP 0 955 714 A2; and U.S. 2002/0067096 A1.\nThe rotor presented in the invention preferably finds application in a brushless DC motor or another permanent magnet synchronous motor. Such motors can be used in a great variety of applications, including spindle motors for disc drives, motor-assisted systems in motor vehicles such as steering and braking systems, electric tools and many other applications.\nThe radial arrangement of the permanent magnets embedded in the rotor core gives rise to the problem of stray flux in the region of the shaft onto which the rotor is mounted. The shaft is usually made from steel and acts as an extra back yoke for the magnetic flux through the rotor core. This gives rise to considerable magnetic stray. This problem can be countered by fitting a sleeve made from a magnetically non-conductive or low-conductive material to the shaft to which the flux guide elements of the rotor core are fixed, between which in turn the permanent magnets are embedded. This construction method is relatively costly and requires extra individual parts.\nThe object of the present invention is to submit a rotor for an electric motor which has embedded magnets and is simple to manufacture but nonetheless prevents the above problem of stray flux being generated in the region of the shaft."}
{"text": "1. Field of the Disclosure\nThe disclosure relates generally to sensory feedback systems and methods for guiding users in virtual/augmented reality environments such as walk-around virtual reality environments and for assisting with preventing collisions with objects in the physical operating space in which the user acts.\n2. General Background\nVarious augmented and/or virtual reality systems and/or environments are known. One current generation of desktop virtual reality (“VR”) experiences is created using head-mounted displays (“HMDs”), which can be tethered to a stationary computer (such as a personal computer (“PC”), laptop, or game console), or self-contained. Such desktop VR experiences generally try to be fully immersive and disconnect the users' senses from their surroundings.\nCollisions with physical objects when using a walk-around virtual reality system are currently solved in certain situations by either having a second person in the operating space (a “chaperone”) guiding the user, and/or by providing physical hints (e.g., by placing a thick carpet on the floor that ends some distance from the adjacent walls).\nIt is desirable to address the current limitations in this art."}
{"text": "Prior art methods of processing text typically incorporate linguistic knowledge that is not resident in the text (e.g., document) being processed. Prior art text summarization methods often rely on a collection of exemplary text that is external to the text being processed to assess the role a word plays in the text being processed. For those methods that rely on a collection of exemplary text, it is difficult, if not impossible, to generate a single collection of exemplary text that can be used to successfully summarize textual documents on widely different topics because a word in one context may have a different meaning in another context. This problem is often overcome in the prior art by generating multiple collections of exemplary text, where each collection is tailored to a specific topic (e.g., scientific, financial). Generating a collection of exemplary text is difficult, time consuming, and prone to error (e.g., biases of those generating the collection).\nU.S. Pat. No. 5,384,703, entitled “METHOD AND APPARATUS FOR SUMMARIZING DOCUMENTS ACCORDING TO THEME”; U.S. Pat. No. 5,638,543, entitled, “METHOD AND APPARATUS FOR AUTOMATIC DOCUMENT SUMMARIZATION”; U.S. Pat. No. 5,708,822, entitled, “METHODS AND APPARATUS FOR THEMATIC PARSING OF DISCOURSE”; U.S. Pat. No. 5,768,580, entitled “METHODS AND APPARATUS FOR DYNAMIC CLASSIFICATION OF DISCOURSE”; U.S. Pat. No. 5,978,820, entitled “TEXT SUMMARIZING METHOD AND SYSTEM”; U.S. Pat. No. 5,924,108, entitled “DOCUMENT SUMMARIZER FOR WORD PROCESSORS”; and U.S. Pat. No. 6,199,034, entitled “METHODS AND APPARATUS FOR DETERMINING THEME FOR DISCOURSE,” each disclose a method of summarizing text that relies on information not contained in the text being summarized. The information not contained in the text comes in various forms such as a stop list, a lexicon, a knowledge catalog that includes static ontologies, a word dictionary, and a pre-compiled list of words and phrases. The present invention does not rely on any of these forms of information. U.S. Pat. Nos. 5,384,703; 5,638,543; 5,708,822; 5,768,580; 5,978,820; 5,924,108; and 6,199,034 are hereby incorporated by reference into the specification of the-present invention."}
{"text": "This application relates to a two-part polyurethane based adhesive useful in bonding porous and nonporous materials. The adhesive is especially useful in bonding glass into window frames, for example windshields into automobiles.\nPolyurethane adhesive compositions typically comprise at least one urethane prepolymer. Adhesives useful for bonding to nonporous substrates, such as glass, to metal, are well known (see U.S. Pat. Nos. 4,374,237 and 4,687,533 incorporated herein by reference). When glass is installed in automobiles on a production line, a one-part moisture curable polyurethane adhesive is preferably used, as the equipment needed for the application of such an adhesive in an automobile assembly plant is less expensive than the equipment needed to apply a two-part adhesive. One part polyurethane adhesives are disclosed in U.S. Pat. Nos. 4,374,237 and 4,687,533.\nIn the automotive after market replacement industry, glass is often bound into automobiles through the use of two-part moisture curable polyurethane adhesives. Two-part polyurethane adhesives are used because they offer rapid initial cure allowing for rapid drive-away times. Two-part polyurethane adhesives are well known, see for example U.S. Pat. No. 4,835,012, incorporated herein by reference, and DE 4,210,277. Despite the relatively rapid drive-away times afforded by two-part polyurethane adhesives, the market demands faster curing adhesives which allow even faster drive away times, for instance about 60 minutes from application and more preferably about 30 minutes from application. If the adhesive cures too rapidly, then the window installer loses the necessary time to install and properly place the glass into the frame before the adhesive becomes too intractable to work with. Working time is defined as the period from application of adhesive until the adhesive becomes too intractable to work with, and is preferably about 10 to 12 minutes.\nFuture regulations in the U.S. will require dual airbags in vehicles. During crashes the airbags inflate and exert additional pressure on the windshield. Federal Transportation Agency regulations require that windshield remain in place at crashes of up to 30 mph (48 KPH). This requires adhesives which have enhanced strength at the drive away time. To meet such a standard, the lap shear strengths of the adhesive should be preferably 150 psi (1033 kPa) at the designated drive away time as determined according to ASTM D-3163.\nTherefore, what is needed is a two-part polyurethane adhesive which facilitates faster drive away time, which meets the strength requirements defined above and which still provides for a reasonable working time to facilitate proper placement of glass in window frames."}
{"text": "This is a 35 U.S.C. 371 application based on PCT/US94/00445, filed Feb. 16, 1994.\nThe invention relates to an apparatus for lengthening bones.\nApparatuses of this kind, which not only lengthen bones but in addition serve to transport bone fragments within part of a limb, to restore the mobility of joints, to correct congenital or acquired axial misalignments of the limbs, to treat fresh or slowly healing bone fractures, to aid the healing of artificial joints and to stiffen joints by acting as an external stabilizer or fixator, are sufficiently well known. Each such apparatus for lengthening bones comprises at least two frame elements at a distance from one another, which at least partly enclose the bone, to receive and fasten both fixation elements and tensioning elements or the like. The fixation elements, which transfer force to the bone and take the form of wires such as Kirschner wires, screws, Schanz's screws etc., transfix the bone about perpendicularly to its long direction or are drilled into and/or through the bone under aseptic, sterile conditions and fixed to the external frame elements. The tensioning elements or the like, which can be extended telescopically and which connect the frame elements to one another and establish their relative positions, are provided to produce controllable and continuous tractive and pressure forces to lengthen bones, treat bone fractures, and mobilize or stabilize joints.\nAll these apparatuses present the disadvantage of decidedly elaborate construction. For example, their frame elements are generally constructed as closed rings provided with a single row of holes to receive and fasten the fixation elements on the one hand and the tensioning elements or the like on the other hand. Apart from the fact that such closed rings are not needed in the treatment of every case, they complicate the manipulation of the whole apparatus, particularly for placement and removal during the surgical operation. Furthermore, such closed rings prevent direct access to the bone concerned. In addition, the arrangement of the fixation elements and tensioning elements or the like at each of the frame elements is largely predetermined by the single row of holes, so that it cannot be altered as desired. As another example, the tensioning elements or the like in these apparatuses are in many cases formed by a threaded sleeve with two threaded rods, one at each of its ends, which can be screwed into the threaded sleeve. By rotation of the threaded sleeve, in many cases both of the threaded rods fastened to the associated frame element are simultaneously lengthened or shortened. With this arrangement, fine adjustment of the distance between the frame elements is very difficult. Moreover, the threaded sleeve in such an apparatus has a fixed, unalterable length, so that a tensioning element or the like intended for a relatively short distance between the associated frame elements cannot be used for a larger distance between the associated frame elements. Such tensioning elements or the like, and hence these apparatuses, cannot be employed in a versatile manner when it is desired to vary the distance between the frame elements, unless the tensioning elements or the like are exchanged for others of different length.\nNot the least result of these disadvantages is that manipulation of these apparatuses is relatively difficult, both during the operation phase and during the subsequent healing phase.\nFinally, in practice such apparatuses have additionally proved to be very heavy, which severely affects acceptance by the patient."}
{"text": "Demand for electric vehicles with hybrid drive has soared worldwide due mainly to a recent sharp increase in fuel prices. Conventional battery management systems for electric vehicles (EVs), however, are designed in an ad hoc way, causing the supply of electric vehicles to fall behind the market demand.\nCost-effective electric vehicles require not only development of high energy density battery cells, but also efficient management of large-scale battery packs, each consisting of a large number of battery cells. In particular, a battery management system (BMS) that monitors and controls battery cells in a pack, must cope with heterogeneous batterycell characteristics. That is, even if characteristics of all battery cells in a battery pack are initially identical, as they are charged and discharged repeatedly, each cell will exhibit different characteristics. A weak cell—hat is (charged and/or) discharged faster than others—is likely to be (over-charged and/or) deep-discharged, i.e., the battery cell continues to be discharged even when its terminal voltage falls below a certain threshold called a cutoff voltage. This weak battery cell can eventually become faulty, and will, if not managed properly, cause the whole pack to be dysfunctional.\nA battery management system should be able to cope with weak/faulty cells in such a way that faulty cells are bypassed to keep the pack operational. Bypassing certain cells inside a pack, however, requires switches by which the connection arrangement of battery cells can be changed. Switches are placed around battery cells, regulating the battery supply power. Furthermore, a reconfigurable battery system may offer a way to alter battery connectivity and dynamically adjust supply power to meet application demands. All of these systems require careful system specification, cost-effective incorporation and control of system components, such as switches and battery cells.\nThere are two main challenges in developing a battery management architecture. First, there is a tradeoff between the minimum number of hardware components to use and maximum reconfigurability in a BMS. Key components therein are switches that allow a battery-cell array to be reconfigurable. The more switches around cells, the more reconfigurable the array becomes, but the costlier. Also, individual components affect directly system reliability. System reliability should be assessed based on the reliability of components and their connections. At the same time, since the cost is the major consideration in realizing a reconfigurable architecture, the components count should be minimized. Second, to maximize both system reconfigurability and reliability, a reconfigurable architecture should be specified with respect to software/hardware components and their inter-relationship. An application (software) may require various battery (hardware) conditions from a BMS. Also, a BMS may request subsystem/local BMSs, if any, for the information on the status of individual battery cells in the case of modular management architecture. Upon receipt of this request, individual local BMSs periodically monitor their battery-cell arrays and reconfigure them, if necessary, in accordance with individual cell characteristics. This interaction between local BMSs also depends upon the underlying hardware system design. A well-designed, combined hardware-software battery management architecture will provide high reliability, cost-effectiveness, and scalability.\nThis section provides background information related to the present disclosure which is not necessarily prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to a pressurized water closet flushing system that minimizes water usage incident to flushing of a toilet yet maximizes the efficiency of effluent transport.\n2. Related Art\nThe herein disclosed pressurized water closet flushing system represents an improvement over the system disclosed in my U.S. Pat. No. 4,233,698 issued Nov. 18, 1980 and is a continuation of my application Ser. No. 08/457,162, filed Jun. 1, 1995 abandoned.\nWater conservation is an environmental problem that has resulted in strict controls being placed on domestic water usage in many areas of the country. Pressurized water closet flushing systems make a significant contribution to water conservation in that they exhibit relatively low water consumption coupled with high effluent transport efficiency. Known systems generally consist of a water supply group, an accumulator vessel, a flush valve and a flush control. The aforesaid components are generally installed internally of a water closet and are energized by water pressure from a fresh water supply system.\nA complete flush cycle consists of a water fill phase and a water discharge phase. In the fill phase of the cycle, supply system pressure forces the water into the accumulator vessel. As the water level rises in the accumulator vessel, air contained therein is compressed. When the pressure of the compressed air in the accumulator vessel equals that of the fresh water supply, flow of water into the accumulator vessel ceases and the system is conditioned for the water discharge phase of operation. When the flush control is actuated, the water discharge phase commences and the compressed air in the accumulator vessel pushes the stored water into the water closet bowl at high velocity, flushing waste therefrom with minimum water consumption.\nThe function of the accumulator vessel is to store both water and potential energy in the form of compressed air. For a given line pressure, the volume of the accumulator vessel determines the maximum discharge energy available.\nThe water closet flushing system of the present invention exhibits a substantial increase in discharge energy over known systems without a corresponding increase in water consumption. Discharge energy is maximized by increasing the volume and therefore the total potential energy of the compressed air charge above the water in the accumulator vessel while water consumption is minimized by positive closing of an improved flush valve.\nThe improved flush valve features a novel balanced piston that divides a flush valve cylinder into upper and lower chambers. As water enters the accumulator vessel, in the fill phase of the flush cycle, compressed air flows from the lower chamber defined by the piston through an air transfer orifice in the piston into an upper chamber of the cylinder above the piston. The fill phase of the cycle continues until a pressure balance is achieved between pressure in the accumulator and line pressure. The flush control is connected to the upper chamber of the flush valve cylinder. When the flush control is opened, the water discharge phase of the flush cycle is initiated by venting the compressed air in the upper chamber of the cylinder creating a pressure differential across the piston allowing the piston and flush valve thereon to move upwardly allowing water in the accumulator vessel to be discharged at high velocity into the water closet bowl.\nClosing or downward movement of the flush valve piston is initially resisted by a partial vacuum that is created within the upper chamber of the cylinder above the piston due to initial downward movement thereof. The flush valve piston remains suspended with the flush valve thereon in the open condition until sufficient air passes from the accumulator to the lower chamber of the flush valve cylinder thence upwardly through the transfer orifice in the piston into the upper chamber of the cylinder to reduce the pressure differential across the piston, allowing the flush valve piston and flush valve to return to the closed position under the bias of the flush valve return spring.\nThe aforesaid operation of the flush valve is insured by an improved air induction system that provides for replenishment of air lost in each flush cycle and lost due to absorption of air into the stored water. Moreover, the air induction system is self-limiting and the accumulator vessel cannot be overcharged with air, e.g., become xe2x80x9cair-logged.xe2x80x9d\nIn accordance with yet another feature, an improved vacuum breaker acts as a safety device that precludes contaminated water in the toilet bowl from being siphoned into the accumulator vessel and then into the fresh water system. In the event that a negative pressure develops within the fresh water supply system resulting in a relatively higher pressure in the accumulator vessel than in the water supply system, water contained within the vessel will flow backwards reducing the pressure within the vessel. When the internal pressure within the vessel falls below ambient pressure, the vacuum breaker valve will open and admit outside air into the vessel breaking the vacuum and precluding waste from the water closet bowl from being pulled into the vessel. Placement of the vacuum breaker above the flush valve piston precludes creation of a vacuum above the piston in the event air in the accumulator vessel is fully expanded which could restrict proper closure of the flush valve.\nAdvantages of the herein disclosed system over prior systems are greater operational reliability, higher efficiency and lower manufacturing cost."}
{"text": "Field of the Invention\nThe present invention relates to a method of manufacturing a sealing device which seals a moving part of a device with a seal lip which is integrally formed in a reinforcing ring.\nDescription of the Conventional Art\nAs a sealing means for a bearing portion which rotatably supports a wheel to a wheel suspension device for a motor vehicle, there has been conventionally known a sealing device, for example, as described in Japanese Unexamined Patent Publication No. 2008-281013.\nFIG. 7 shows a sealing device which is the same kind as that described in Japanese Unexamined Patent Publication No. 2008-281013 in a used state, reference numeral 110 denotes an outer ring of a bearing which rotatably supports a wheel to a wheel suspension device for a motor vehicle, and reference numeral 120 denotes a hub in which a center shaft portion 122 thereof is inserted to the outer ring 110 and the wheel (not shown) is attached to an attachment flange 121. A plurality of spherical bodies (not shown) arranged at predetermined intervals in a circumferential direction are retained to a bearing internal space S in a lower left side in the drawing so as to freely roll on a track (not shown) which is formed in an outer peripheral surface of the center shaft portion 122 and an inner peripheral surface of the outer ring 110, and a bearing unit is constructed by the center shaft portion 122 of the hub 120, the outer ring 110 and the spherical bodies.\nA sealing device 100 is interposed between an inner peripheral surface of an outer end portion of the outer ring 110 and the outer peripheral surface of the center shaft portion 122 in the hub 120. The sealing device 100 is structured such as to prevent a foreign material such as rain water, muddy water or dusts from intruding into the bearing internal space S from an external portion A, and prevent grease filled in the bearing internal space S from leaking to the external portion A.\nIn detail, the sealing device 100 is constructed by a reinforcing ring 101 made of metal, a labyrinth lip 102, anti-dust lips 103 and 104, and an anti-grease lip 105 as shown in FIG. 8. The labyrinth lip 102, the anti-dust lips 103 and 104 and the anti-grease lip 105 are made of a rubber-like elastic material (a rubber material or a synthetic resin material having a rubber-like elasticity), and is integrally vulcanization bonded to the reinforcing ring 101 via a coating layer 106 which is made of the same rubber-like elastic material.\nFurther, the labyrinth lip 102 in the outermost diameter side is moved close to and faced to an end surface of the attachment flange 121 of the hub 120 in a non-contact manner, the labyrinth lip 102 in the outermost diameter side being positioned in a back surface side of an outer diameter flange portion 101b which is developed like a disk to an outer peripheral side of the outer ring 110 from an end portion of a fitting tube portion 101a which is pressed into and fitted to an inner peripheral surface of the outer ring 110. As a result, the reinforcing ring 101 serves as a labyrinth seal."}
{"text": "The present disclosure pertains to an apparatus and a method for fiber delivery of high power laser beams.\nOptical fibers can be used to transmit a laser beam from a laser source to a desired location. The use of an optical fiber to transmit a laser beam from a laser source to a desired location is a significant enabler in a number of laser applications because the optical fiber offers a flexible transmission path that involves no free-space optics and which can be re-routed in real time. For example, industrial lasers operating at wavelengths that can be transmitted through silica fibers almost universally use such optical fibers to move the laser beam from the laser source to the workpiece. Often, the output end of the optical fiber is attached to a robotic arm that directs the output beam to the workpiece. Another example is in medical applications, where an optical fiber is used to transmit a laser beam through a blood vessel to a specific location where the focused beam is used to accomplish a desired medical effect. Another example appears in solid-state laser designs, where radiation from pump diodes is commonly delivered to active region of a solid-state laser via fibers or bundle of fibers. Yet another example is in defense systems, where guiding a laser beam onto the inner gimbal of a beam director is desired without having to propagate through a complex sequence of optical elements that move and guide the laser beam.\nWhile the first three examples above are already implemented in their respective application areas, the use of optical fiber delivery of laser beams in defense systems has been limited in the past by two challenges. First, the optical fiber in conventional systems is required to have a sufficiently large core area so that it can accommodate the required average or peak power without suffering from optical damage or degradation of the laser beam by nonlinear optical processes such as stimulated Brillouin scattering (SBS) or stimulated Raman scattering (SRS). Second, if a sufficiently large core is developed, the fiber must be designed such that the input beam quality, which for many military applications must be nearly diffraction-limited, is not degraded by propagation through the fiber.\nOptical fibers are widely used for laser beam delivery in the industrial laser market. These optical fibers are capable of delivering laser beams with very high powers of interest. However, these conventional optical fibers are designed to transport radiation of highly multimode industrial lasers. These conventional optical fibers necessarily have core diameters as large as 0.5 to 1 mm and a large numerical aperture (NA) 0.1-0.2, such that the fibers are also highly multi-mode. Presently, there is no practical way to maintain high beam quality when using such conventional optical fibers for beam delivery. Due to the multimode aspect of these conventional optical fibers, bending of the optical fibers results in a strong mode coupling to the higher order transverse electromagnetic modes that are guided along with the fundamental and other lowest order transverse electromagnetic modes. If the lowest-order transverse electromagnetic mode is launched in these conventional optical fibers, the lowest-order mode will lose most of its power as it feeds higher-order transverse electromagnetic modes. As a result, even if the input beam is nearly diffraction-limited and the output beam suffers only minor power loss, the output beam quality is typically greater than 50-100 times diffraction limited (XDL). Therefore, these conventional optical fibers meet power-delivery requirement, but not beam-quality delivery requirements.\nAnother problem with typical beam-delivery fibers for high power lasers is that as the core diameter increases to accommodate the increasingly high laser power, the fiber becomes less flexible. At some point the basic purpose of the fiber, mechanical flexibility, becomes severely restricted due to the large fiber diameter.\nTo overcome the above deficiency with high beam quality transportation, a large mode area (LMA) optical fiber design can be implemented. LMA optical fibers can guide a few higher-order transverse electromagnetic modes while still maintaining beam quality at or better than approximately 1.3 times diffraction-limited (XDL). A LMA optical fiber differs from the standard large-core delivery fibers by having a relatively small core diameter, between about 20 μm and about 30 μm for signal wavelengths of about and a reduced NA of approximately 0.06. LMA fibers must be properly coiled to maintain good beam quality. At a prescribed core diameter, coiling induces bending losses to all transverse electromagnetic modes, but the higher-order modes all suffer much greater coil-induced loss than the desired lowest-order transverse mode. Hence, coiling results in increased radiation loss for higher-order transverse electromagnetic modes, which are stripped out of the core, thereby “cleaning up” the laser beam and yielding a beam having mostly the fundamental transverse electromagnetic mode. With an optimized coil radius, the loss for the fundamental transverse electromagnetic mode remains at a low and tolerable level.\nA 10 kW fiber laser that generates a laser beam to be delivered to a beam director for use in industrial applications is often based on an LMA fiber. A conventional LMA beam delivery fiber can be considered to meet the requirements of beam delivery. Indeed, since the laser beam is generated by the LMA fiber in the first place, the LMA fiber can obviously accommodate the laser power. However, there are other considerations that need to be taken into account. Among these considerations is the onset of stimulated Raman scattering (SRS) or stimulated Brillouin scattering (SBS) as the beam-delivery fiber length approaches a few meters, as is required for typical applications.\nWhat is needed is a method and apparatus that cure the deficiencies noted above in conventional optical fibers used for delivery of high power laser beams, and that provide an optical fiber taking into account the above considerations, including SRS and SBS."}
{"text": "Demand for bandwidth is driving the expansion of optical transmission systems into homes and businesses of all sizes. Single wavelength fiber optic systems can support substantial data rates. However, services such as HDTV, on-demand TV programming, internet telephony, and telepresence are bandwidth intensive beyond the capabilities of many traditional networks. The present invention relates to integrated wavelength selectable photodiodes and their application to Fiber-To-The-X (FTTX) services. Fiber-To-The-X services refer to the extension of optical data transport into areas traditionally served by electrical communications systems, such as homes and small and medium sized businesses. Examples of FTTX systems are Fiber-To-The Home (FTTH), Fiber-To-The Curb (FTTC) and Fiber-To-The-Premises. FTTX architectures are also used for some highly secure optical communications links, such as radar tower interfaces."}
{"text": "The digital laser range finder, because of its versatility and accuracy, is widely used today for computing ranges to targets, landmarks and various objects. Unfortunately, the dangerous nature of lasers limits the use thereof to only those personnel that have been well trained in their operation. Further, availability and cost may restrict the use of laser range finders in many operations. Such safety concerns, laser availability, and cost of the laser range finder often make it difficult to develop systems that require the integration of the digital laser range finder and its output.\nWhile it would appear that a proper solution to this problem would be to simulate the serial data output signal of a laser range finder, the implementation of this solution has proved rather difficult. For example, use of a programmable pulse generator or personal computer to generate serial signals indicative of a laser range finder is suitable for lab conditions. However, at some point it becomes necessary to field test an entire system of which the laser range finder is an integral part. For such field tests, the typical \"lab\" type devices are too cumbersome, present power supply problems, and are not designed to withstand field test environments in which an entire integrated system needs to be tested. Furthermore, these devices cost thousands of dollars and are typically \"over qualified\" for the task at hand."}
{"text": "The present invention relates to fire resistive foamed polystyrene herein referred to as xe2x80x9cStyrofoamxe2x80x9d products and a method of their production.\nOne of typical examples of conventional methods of producing a Styrofoam product calls for producing partially formed beads by conducting preliminary forming of numerous styrene beads that have been produced by forming polystyrene into the shape of beads, and then, after leaving these pre-expanded beads to age, expanding them properly.\nAs conventional Styrofoam products having thus produced are not only light in mass but also non-hygroscopic, superbly heat insulating, easy to handle and very inexpensive, they enjoy a tremendous demand and are widely used in the field of ordinary houses, building materials, miscellaneous goods or the like.\nConventional Styrofoam products, however, have drawbacks such that they have little ability to withstand fire and, when exposed to fire, are easy to become contracted and deformed so that their entire shape collapses instantly.\nIn order to solve the above problem, the present invention relates to Styrofoam products that have superior fire resistance and are capable of limiting degree of contraction and deformation caused by fire. The present invention also relates to a method of producing such Styrofoam products.\nA Styrofoam product according to the present invention includes one or more formed styrene beads and a coating film or coating films that contain boric acid-based inorganic materials and are formed on the surface(s) of the styrene bead(s). As a coating film or coating films containing boric acid-based inorganic materials are thus formed on the surface(s) of the formed styrene bead(s), the coating film(s) suppress burning of the styrene bead(s). In case a great quantity of heat is applied, the vitreous coat into which each coating film has chemically changed prevents the styrene bead from burning deeper. Therefore, the invention increases the heat resisting ability of the Styrofoam product and prevents its thermal contraction and deformation.\nA Styrofoam product according to another feature of the present invention includes a plurality of styrene beads that have been expanded and solidly bonded to one another, and coating films that contain boric acid-based inorganic materials and are respectively formed on the surfaces of the styrene beads. As coating films containing boric acid-based inorganic materials are thus formed on the surfaces of the styrene beads, which have been expanded and solidly bonded to one another, the coating films suppress burning of the styrene beads. In case a great quantity of heat is applied, the vitreous coats into which the coating films have chemically changed prevent the styrene beads from burning deeper. The invention thus increases the heat resisting ability of the Styrofoam product and prevents its thermal contraction and deformation.\nA method of producing a Styrofoam product according to the invention calls for producing a beads mixture by mixing styrene beads with a boric acid-based inorganic material, producing partially foamed beads by pre-expanding the beads mixture, producing additive-containing beads by adding boric acid-based inorganic materials to the partially foamed beads, and conducting a main expansion treatment to the additive-containing beads. As a result of the above method, coating films containing boric acid-based inorganic materials are formed on the surfaces of the foamed styrene beads, and the coating films increase the fire resistive ability of the Styrofoam product and prevent its thermal contraction and deformation.\nA method of producing a Styrofoam product according to another feature of the invention calls for producing a beads mixture by mixing styrene beads with a boric acid-based inorganic material, producing partially foamed beads by pre-expanding the beads mixture, producing additive-containing beads by adding boric acid-based inorganic materials and a thermosetting resin to the partially expanding beads, and conducting a main expansion treatment to the additive-containing beads. As a result of the above method, coating films containing boric acid-based inorganic materials and a thermosetting resin are formed on the surfaces of the formed styrene beads. Therefore, because of synergistic effect between the boric acid-based inorganic materials and the thermosetting resin, the coating films increase the fire resistance of the Styrofoam product and reliably prevent its thermal contraction and deformation.\nA method of producing a Styrofoam product according to yet another feature of the invention calls for producing a beads mixture by mixing styrene beads with a boric acid-based inorganic material, producing partially formed beads by pre-forming the beads mixture, producing additive-containing beads by adding a fibrous material, as well as boric acid-based inorganic materials and a thermosetting resin, to the partially formed beads, and conducting a main expansion treatment to the additive-containing beads. As a result of the above method, coating films that contain boric acid-based inorganic materials and a thermosetting resin and are formed on the surfaces of the formed styrene beads increase the fire resistance of the Styrofoam product and reliably prevent its thermal contraction and deformation. Furthermore, the fibrous characteristics of the fibrous material increases the strength of the product.\nA method of producing a Styrofoam product according to yet another feature of the invention calls for producing a beads mixture by mixing styrene beads with a boric acid-based inorganic material, producing pre-formed beads by preliminary forming the beads mixture, producing additive-containing beads by adding either one of or both an amino-based resin and a polyamide resin, as well as boric acid-based inorganic materials, a thermosetting resin and a fibrous material, to the pre-formed beads, and conducting a main expansion treatment to the additive-containing beads. As a result of the above method, coating films that contain boric acid-based inorganic materials and a thermosetting resin and are formed on the surfaces of the formed styrene beads increase the fire resistance of the Styrofoam product and reliably prevent its thermal contraction and deformation. Furthermore, the fibrous characteristics of the fibrous material increases the strength of the product, and the amino-based resin and/or the polyamide resin make the product tough."}
{"text": "Cloud computing is an Internet-based computing concept whereby shared resources, software and information are provided to computers and other devices on-demand, like a public utility.\nThe term “cloud” is used as a metaphor for a network, based on the cloud drawing used to represent the telephone network, and later to depict the Internet in computer network diagrams as an abstraction of the underlying infrastructure it represents. Typical cloud computing providers deliver common business applications online which are accessed from another web service or software, like a web browser, while the software and data are stored on servers.\nIn general, cloud computing customers do not own the physical infrastructure. Instead customers can avoid capital expenditure by renting usage from a third-party provider. They consume resources as a service and pay only for resources that they use. Many cloud-computing offerings employ the utility computing model, which is analogous to how traditional utility services (such as electricity) are consumed, whereas others bill on a subscription basis. Sharing “perishable and intangible” computing power among multiple tenants offer the promise of improving utilization rates, as servers are not unnecessarily left idle (which can reduce costs significantly while increasing the speed of application development)."}
{"text": "A trunked communication system is a radio mobile communication system in which multiple users share multiple channels via dynamic channel assignment. The system generally consists of mobile terminals and a dispatcher station, and has functions of dispatch, group call, priority call, etc. In a conventional trunked communication system, the dispatcher station is mainly responsible for dispatch and smart distribution of mobile terminals. When the mobile terminal is not able to access an incoming call, for example, when the mobile terminal is busy or the battery runs out, the dispatcher station forwards the call to the mobile terminal as needed.\nAfter registering in the dispatcher station, the mobile terminal applies to the dispatcher station for call diversion; in response to an instruction from the dispatcher station for obtaining target address, the mobile terminal sends the target address to the dispatcher station; and the dispatcher station stores the received target address. A caller terminal requests to the dispatcher station calling the mobile terminal, and if the mobile terminal is not able to access the incoming call, for example, when the mobile terminal is busy or battery runs out, the dispatcher station sends the target address to the caller terminal, the caller terminal initiates a call to the target terminal, and the call diversion of mobile terminal is accomplished.\nWhen using the above method to implement call diversion, those skilled in the art find the following drawbacks.\nThe mobile terminal sends the target address to the dispatcher station to establish call diversion setting, and when performing the call diversion, the dispatcher station sends the target address to the caller terminal to implement call diversion. The mobile terminal is passive in the process of call dispatch and distribution, and is controlled by the dispatcher station. In a conventional communications system without the dispatcher station, the mobile terminal can only implement basic call services and information transmission function, and cannot implement call diversion."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrolytic capacitor comprising an anode of a valve metal such as aluminum or tantalum, and a cathode of a conductive polymer layer, and also relates to a method of producing such an electrolytic capacitor.\n2. Prior Art\nA conventional electrolytic capacitor have been composed of an anode formed of the valve metal, a dielectric layer formed of a film of the valve metal oxide, and a cathode formed of a liquid electrolyte or inorganic solid electrolyte, with leads connected to the anode and the cathode, respectively, being packaged to make a product.\nThe anode is made of aluminum or tantalum as the valve metal, whose oxide metal is used for the dielectric layer. In the case of aluminum electrolytic capacitor, the surface of an aluminum foil is coarsened and anodized to form the oxide film, thereby obtaining an aluminum anode element. In the case of a tantalum electrolytic capacitor, a compact made of powdered tantalum is fired in a oxidizing atmosphere, to obtain a tantalum anode element.\nConventionally, as a cathode, the organic solvent including an organic acid or the like have been used in the case of the aluminum electrolytic capacitor. The solid electrolyte was, for example, manganese dioxide in the case of a tantalum electrolytic capacitor.\nIn these days, electronic components have been requested to have excellent high-frequency response property, particularly, to be adaptable to digital circuits, so that electrolytic capacitors are also required to have excellent high-frequency response and lower inner resistance. In order to meet these market needs, some conductive polymers having a high conductivity have been examined and developed for use in the solid electrolytic cathode of such electrolytic capacitors.\nIn the electrolytic capacitor wherein the solid electrolyte is formed of a conductive polymer, the conductive polymer layer is formed on the anode element to make direct contact to the entire dielectric layer of the oxide film on the valve metal.\nThe electrolytic capacitor using such a conductive polymer layer is known, for example, in the Japanese Patent Publication JP-A 6-168855.\nThe chemical polymerization technique have been known, which forms the conductive polymer on an anode elements into a cathode. In this technique, the anode element is immersed in a monomer solution, and the monomer is polymerized at the presence of an oxidizing agent as a polymerizing initiator. As a result, a conductive polymer layer is formed on the oxide film on the surface of the anode.\nThe electrolytic polymerization technique have also been known for polymerizing, wherein polymerization is performed by anodizing a conductive monomer on the surface of a solid film to have polymerized as the electrolyzing electrode. According to this technique, a first conductive layer, having slight conductivity, must be formed on a dielectric oxide film. This first conductive layer on the anode element is used a start point for further polymerizing electrolytically. For this first conductive layer is used a manganese dioxide film previously attached to the oxide film, or a very thin conductive polymer layer previously formed by the above chemical polymerization technique.\nHowever, the formation of the conductive polymer layer by the conventional techniques caused some problems. It was difficult to grow a conductive polymer layer having a uniform and sufficient thickness on the entire surface of the dielectric layer of the oxide film. For this reason, the electrolytic capacitor obtained as the final product could not show capacitance high enough to be estimated from the anode film design. Furthermore, in the prior art chemical polymerization technique, the polymerizing operations had to be repeated several tens of times to increase the thickness of the grown polymer, in order to obtain the sufficiently high capacitance.\nIn the case of using the conventional electrolytic polymerization technique, the first conductive layer is formed of a manganese dioxide or a very thin conductive polymer layer formed by the chemical polymerization technique. Therefore, electric current supplied for electrolysis flows through the first conductive layer having a relatively high electric resistance. For this reason, a relatively long polymerization period of time was required to form the conductive polymer layer on the entire surface of the dielectric layer.\nConventionally, to electrolytically form a conductive polymer layer on both sides of an anode valve metal foil, polymerization had carried out, starting from a part of one side of the anode valve metal foil, and then advancing to the opposite side through roundabout routes. In this conventional method, when a conductive polymer layer is formed by electrolytic polymerization on both sides of an anode foil having a relatively large area, it takes a long time for polymerization. In addition, the thickness of the conductive polymer layer at the starting point for polymerizing differs from that at the ending point. This difference in thickness of the conductive polymer layer is caused by the total polymerizing period of time different at the start point and the end point. Therefore, to obtain a conductive polymer layer having a more relatively uniform thickness, the speed of polymerization is required to be made slower, then reducing the difference in polymerization time. This results in extending the operating time for polymerization.\nFor these reasons, the above technique, not to have formed a conductive polymer layer uniformly on a wide anode foil, was limited to only on a relatively narrow anode foil which is suited for the shape of the final product, just as in the example of the Japanese Patent Publication JP-A 63-239917.\nIn order to form a necessary and sufficient conductive polymer layer for the final product, the processes for capacitor production have been complicated and taken a long time, then, requiring high production cost."}
{"text": "A dual-type wireless communication terminal has been known as a wireless communication terminal capable of performing communication by switching between two communication systems.\nPatent Document 1: JP-A-H9-172675\nOf the above dual-type wireless communication terminals, a wireless communication terminal, which can measure a wireless communication status of one communication method during communication using another communication method and monitor an incoming call, is especially called a hybrid type.\nIn relation to the hybrid-type wireless communication terminal, a hybrid-type wireless communication terminal using a cdma2000 1× system which is mainly designed for voice communication and a 1×EVDO system which is specifically designed for data communication suspends the data communication of the 1×EVDO system at predetermined intervals (e.g., 5.12 seconds) and then switches an antenna and a radio section to the cdma2000 1× to perform system monitoring with the cdma2000 1×, in order to monitor an incoming call of the cdma2000 1× during the data communication with the 1×EVDO. When the system monitoring has been completed, the antenna and the radio section are again switched to the 1×EVDO, thereby resuming the data communication with the 1×EVDO. In the case that the wireless communication terminal is located in an area where is in the vicinity of a boundary between service areas by a plurality of base stations as cdma2000 1× system (the case that quantities of signals received from the plurality of base stations (C/I values) contend with each other), the cdma2000 1×, thereby performing the system monitoring with the cdma2000 1×. At this time, when the idle handoff of the cdma2000 1× system has been detected and the idle handoff is continuously detected even after the idle handoff processing, the antenna and the radio section of the wireless communication terminal are still occupied with the cdma2000 1× system. Therefore, after a given period of time has elapsed while the occupation is maintained, the wireless communication terminal as the 1×EVDO system determines that radio wave is lost (see FIG. 4). For this reason, during data communication of the 1×EVDO, the data communication of the 1×EVDO system is disconnected as a result of repetition of the idle handoff in the cdma2000 1× system, although the state of radio wave of the 1×EVDO system is good."}
{"text": "With the advance and development of sophisticated communication techniques, especially those involving considerably high data rates as may be encountered in high speed digital communication networks, particularly high speed optical digital communication systems, spare capacity or backup facilities are usually made a part of the network, in order to provide the network with the capability of maximizing the integrity of communication paths through it in the event of an outage or failure of an active communication link. A common solution to the problem of channel faults or data degradation is the provision of one or more protection channels that are normally maintained in a quiescent or standby condition; in response to instructions from channel monitoring equipment a protection channel may be substituted in place of an active channel that has been observed to have developed a fault or an adverse thruput condition that has led or could lead to a loss of data communication over the system.\nConcomitant with the provision of such protection channel equipment is the need to ensure that the substitution of one channel for another is effected rapidly and without adversely affecting the data being transmitted over the channel. For low data rates the time window within which the channel substitution may occur permits a considerably wider error margin and does not necessarily require an observation of the clock signals themselves during the switching process. However, with the very high data rates of present day communication systems (e.g. 200-300 Mb/s) the degree of freedom allowed in switching between channels is significantly restricted. Still, as the demand for information transmission continuity is of paramount importance, some scheme must be provided to accurately control, within an extremely fine error window, the time at which the substitution of one channel for another is effected.\nOne such communication system in which the need arises is described in copending patent application Ser. No. 149,291, filed May 12, 1980, entitled Multi-channel, Repeatered, Fiber Optic Communication Network Having Fault Isolation System by P. Casper et al, and assigned to the assignee of the present application. The communication network described in that application employs a plurality of normally active optical communication channels each of which carries high speed (301 Mb/s) digital data, such as digital telephone traffic. One or more protection channels are provided to carry the digital traffic in the event of a detected problem with the normally active data channels. Advantageously, as described in that application, a channel switching system of the type to be explained in detail below is incorporated in the network in order that what is termed \"hitless\" switching between the active and protection channels may be effected. By \"hitless\" switching is meant the substitution of one data channel for another without the addition of or deletion of data bits from the link. Of course, this \"hitless\" or synchronized switching technique has applicability to data communications other than that described in the above cited application; and it is to be understood, therefore, that reference to such a network is not to be considered limitative of the invention but merely to illustrate an exemplary environment in which the present invention may be incorporated."}
{"text": "In general, the microarrays of single-stranded oligonucleotides (usually DNA or PNA) immobilized on the surface of glass or other solid surfaces are used for various purposes: identification or detection of infectious agents such as virus, bacteria, and other forms of microorganism, genotyping, disease causing genes, mutation of genes, expression of genes, etc. In some instances, the technology has also been adapted for diagnosis at hospitals and clinics. In other instances, various technologies have been consolidated and optimized for convenient use and or reliable decisions in connection with clinical diagnosis.\nThe use of DNA chip often requires long processes such as (a) extraction of DNA or RNA from biological specimens such as tissues, cells, blood, serum, and biological fluids, (b) amplification of target regions of genes by polymerase chain reaction (PCR), or conversion of RNA to cDNA and subsequent amplification of cDNA, (c) labeling of amplified DNA with fluorescent tags, often during PCR, (d) hybridization of the amplified DNA with the oligonucleotides (often called capture probes) immobilized on the surface of a solid substrate such as glass, and (e) scanning of the slide in a detector such as a laser scanner.\nWith respect to extraction, DNA or RNA is obtained from biological specimens by well established protocols (Molecular Cloning: A laboratory manual by Sambrook and Russel, 3rd Edition, Cold Spring Harbor Laboratory Press, 2001). Typically cells are lysed in a buffer containing SDS and proteinase K for degradation of proteins, the lysate is extracted with phenol/chloroform, and the DNA or RNA in the phenol layer is precipitated with ethanol.\nSeveral companies have introduced more simplified protocols for high throughput purification of DNA and RNA from biological specimens. In general, these procedures take advantage of the binding property of DNA and RNA to silica powder, silica fiber or silica membrane (GE Heatlth, Qiagen, etc.). These procedures, although can be accomplished in 30 minutes to one hour, requires multiple pipetting, and many manipulations.\nAs an example, one protocol is briefly described. Biological specimens such as cells are suspended in 200 μl PBS (phosphate buffered saline) and centrifuged (300 g for 5 min). To the pelleted cells 200 μl lysis buffer and 20 μl proteinase K is added. The mixture is incubated at 70° C. for 10 min. Ethanol (200 μl) is added and the tube is vortexed. The lysed cells are transferred into the Mini spin column (with a silica membrane bottom) placed into a 2 ml collection tube and centrifuged at 6000×g for 1 minute and the collection tube is discarded. The Mini-spin column is transferred to another collection tube, 500 μl of wash buffer is added to the Mini spin column and centrifuged. The collection tube is discarded and the Mini-spin column is transferred to another collection tube. A second wash buffer (500 μl) is added and the tube is centrifuged. The Minispin column is transferred to another collection tube, and 200 μl of elution buffer is added to the mini-spin column followed by centrifugation. The collection tube is saved. The Minispin column is then transferred to another collection tube and second elution buffer is added to the Mini-spin column and centrifuged again. The collection tube is saved. The DNA in two collection tubes are combined. As seen from this example, extraction of DNA from cells requires about 10 pipetting steps and several centrifugation steps and many manipulations such as transfer of Mini-spin column to different collection tubes, etc. Purification of RNA also requires a similar process.\nSimilar protocols are adopted for extraction of nucleic acids from multiple samples in an automated format.\nOften, the target DNA or RNA sequences for diagnosis are present in biological samples in a very minute amount, and there is need for amplification of target DNA or RNA sequences before nucleic acid diagnosis. One method for amplification of target DNA sequences is polymerase chain reaction (PCR). Generally, the PCR process involves denaturation of DNA (strand separation) at a high temperature (95° C.), annealing of short primers (usually 15-20 nucleotides) that are complementary to either ends of a DNA region to be amplified, and chain elongation from the annealed primers in the presence of thermostable DNA polymerase and four deoxynucleotide triphosphates. These processes require specific temperatures: for example, 90° C. or above for denaturation, 55-65° C. for annealing, and 72° C. for chain elongation. Denaturation, annealing and chain elongation are repeated continuously until desired amount of DNA is amplified, typically 20-30 cycles. Therefore, PCR requires a programmable thermocycler.\nFor amplification of a RNA sequence, the RNA is first converted to double-stranded cDNA in the presence of reverse transcriptase and DNA polymerase, DNA primers and four deoxynucleotide triphosphates. Subsequently, cDNA is amplified by PCR as described above. The process of cDNA synthesis and subsequent PCR may be carried out in one reaction mixture.\nAfter PCR or cDNA synthesis and subsequent PCR, the amplified product is purified to remove the remaining primers and deoxynucleotide triphosphates. The PCR products are often identified by gel electrophoresis and staining of the amplified DNA with ethidium bromide or CYBR GREEN® dye. [spelled Sybr Green in claim 8]\nFor use of the amplified DNA in diagnosis by DNA chip, the amplified DNA is labeled. For example, a fluorescently labeled deoxynucleotide (Cy3 or Cy5-deoxy nucleotide triphosphate, ALEXA FLUOR®-deoxynucleotide triphosphate) is added to PCR reaction mixture. PCR primers that are labeled with fluorescent group or biotin can also be used for labeling amplified DNA. The biotin-labeled DNA can be detected by interaction with labeled streptavidin. The fluorescently labeled DNA can be detected by laser scanner.\nAfter PCR, the labeled PCR product is purified by use of various protocols. One of which can be the protocol described above by use of silica membrane for extraction of DNA from cells and tissues.\nThe labeled PCR product is added to hybridization buffer. The mixture is heated to 95° C. to convert the labeled amplified DNA to single-stranded form and quickly chilled on ice and transferred on to a chamber assembled on the surface of a chip slide (usually a glass slide) that has immobilized single-stranded oligonucleotide (15-70 nucleotides) capture probes. The slide is then incubated at a temperature for annealing (such as from 40-65° C.). After hybridization, the chamber is disassembled and the slide is washed once with 1×SSC+0.1% SDS, once with 0.1×SSC+0.1% SDS and once with 1×SSC, all as non-limiting examples.\nA DNA chip is typically a slide on which capture probes (short single-stranded DNA) are immobilized in a high density format. In many cases, a glass slide is coated with chemicals to attach the capture probes to the glass. The chemicals should not only attach the DNA to the glass surface but also minimize non-specific binding and signal noise. For example, the chemicals on a chip slide contain silanated, silylated, or poly-L-lysine as a source of an amine or an aldehyde group for attachment of DNA. Recently, dendrons have been developed for use in a biochip. When a chip is coated with dendrons, one can control spacing between capture probes and this also allows reduction of steric hinderance and concomitant increase in sensitivity. Korean patent 10-0383080-0000, and published U.S. Patent Application 2005/0037413 describe dendrons that provide controlled spacing as well as density of amines on dendron.\nSingle-stranded oligonucleotides whose length can vary from 15 to 70 nucleotides are often used as capture probes that will hybridize to the complementary strand sequences of target genes, DNA or RNA. PNA may also be used as capture probes. The oligonucleotides will have either amine or SH group at 3′ or 5′ ends, usually at 3′ end so that the oligonucleotides can be crosslinked to the surface of chip slide.\nWhen DNA chip slides are used, the sequence of steps is to amplify the DNA (labeling the DNA with fluorescent group at the same time) and then purify the DNA. The amplified DNA is denatured by heating and added to the surface of capture probes for annealing. After hybridization, the slide is washed and the fluorescent double-stranded DNA is detected by scanning and the data are analyzed.\nAs evident from above, the whole process from extraction of nucleic acids, gene amplification, and subsequent hybridization of the amplified DNA with the capture probes on chip slide requires multiple processes such as many repeated transfer (pipetting), mixing of solutions, centrifugations, etc. Even when automated machine is developed for extraction of DNA and RNA, there is still need for coupling of this technology with PCR and hybridization with capture probes on chip surface.\nAdditionally, there is a requirement for considerable time for gene amplification and purification (as much as one day). Also, the cost of chemicals and disposable items is high. There is also loss of sample during purification of reaction products. Therefore, there is need for streamlining the cumbersome processes involved in use of gene chips to reduce losses and variation of the results.\nThere have been attempt to carry out PCR in a single step from biological specimens. For example, cells are directly added to PCR reaction mixture. Several companies have introduced one-step PCR mixtures that allow amplification of DNA from whole cells or blood. These mixtures contain detergent to facilitate lysis of cells, proteinase K for digestion of proteins, and often proprietary agents. The efficiency of PCR is variable.\nReverse transcription of RNA to synthesize cDNA with subsequent PCR can also be carried out in one step with whole cells. However, the efficiency again varies depending on the type of biological specimens. For example, blood and serum contains inhibitor of reverse transcriptase.\nThere have also been attempts to couple DNA and RNA extraction by an automated process and PCR. Several companies introduced combined procedures for automated extraction of DNA or RNA and automated PCR. Although these automated processes satisfy diagnosis of certain diseases as well as quantitation of target nucleic acid sequences, such as those of infectious biological agents, it is difficult to use PCR for identification of different variants of a virus (genotypes), changes in sequences in response to drug-resistance, and the mutations of target genes that cause diseases.\nA DNA chip offers many advantages over diagnosis by PCR or by immunologic assays. With a DNA chip, it is possible to identify, all at once, many different sequences, changes of sequences by mutation, different genotypes of viruses, and expression of different genes. However, there is clear need for simplification and integration of all necessary protocols for extraction of nucleic acids, gene amplification and hybridization on chip surface in order for the DNA chip technology to be used at hospitals and clinics with increased throughput and reduction in manual labor. Current technology requires highly trained technicians, long processes, high cost and time. Even when automated systems are available, it requires at least three different systems, such as automated extraction of nucleic acids, gene amplification, and hybridization on DNA chip.\nThere is also an increasing need for diagnosis of disease-causing agents such as viruses and the changes in gene sequences that cause formation of tumors for personalized drug therapy or vaccination. For example, GLEEVEC® (imatinib) is highly effective for treatment of some cases of chronic myeloid leukaemia (CML) and also gastrointestinal stromal tumour (GIST). Gleevec specifically blocks ATP-binding sites on specific mutant tyrosine kinases in the responsive cancers. Therefore, it is necessary to detect the mutant tyrosine kinase genes before treatment of the cancer. Similarly, various oncogenic DNA and RNA viruses are known to induce tumor formation in humans. Prevention of these cancers, such as cervical cancer, requires early detection and treatment of pre-cancerous disease.\nPapiloma viruses are a diverse group of DNA-based viruses that infect the skin and mucous membranes of humans and a variety of animals. Over 100 different human papilloma virus (HPV) types have been identified. Some HPV types may cause warts while others may cause a subclinical infection resulting in precancerous lesions. All HPVs are transmitted by skin-to-skin contact.\nA group of about 30-40 HPVs is typically transmitted through sexual contact and infect the anogenital region. Some sexually transmitted HPVs (types 6 and 11) may cause genital warts.\nPersistent infection with a subset of about 13 so-called “high risk” sexually transmitted HPVs, including types 16, 18, 31, 33, 35, 39, 45, 51, 52, 56, 58, 59, and 68 (different from the ones that cause warts) may lead to the development of cervical intraepithelial neoplasis (CIN), vulvar intraepithelial neoplasia (VIN), penile intraepithelial neoplasia (PIN), and/or anal intraepithelial neoplasia (AIN). These are precancerous lesions that can progress to invasive cancer. HPV infection is a necessary factor in the development of nearly all cervical cancer (Walboomers, J M, Jacobs M V, Manos M M et al (1999) “Human papillomavirus is a necessary cause of invasive cervical cancer worldwide” J. Pathol. 189:12-9).\nAccording to the Center for Disease Control (CDC), by the age of 50 more than 80% of the American women will have contracted at least one strain of genital HPV. All women are encouraged to get a yearly pap smear solely to detect cellular abnormalities caused by HPV. About 14,000 women in the United States are diagnosed with cervical cancer disease each year, and more than 3,900 women die in the United States each year from this disease.\nAmong the high risk HPV strains, type 16 and 18 are together responsible for over 65% of cervical cancer cases.\nAn HPV test detects certain types of HPVs, depending of the test. A method for detecting the DNA of high-risk HPVs has recently been added to the range of clinical options for cervical cancer screening. In March 2003, the US FDA approved a “hybrid-capture” test (see U.S. Pat. No. 6,228,578 B1), marketed by Digene, as a primary screening tool for detecting high-risk HPV infections that may lead to cervical cancer. This test was also approved for use as an adjunct to Pap testing, and may be ordered in response to abnormal Pap smear results.\nThe principle of Digene test is briefly summarized as follows. An exfoliated cervical cell sample is collected in a collection device and nucleic acids are released therefrom. A diluent containing multiple RNA probes to different HPV types is added. The mixture is incubated to allow hybridization of the RNA probes to the HPV sequences. The mixture is then combined with an immobilized antibody for DNA-RNA duplex, and complexes allowed to form. RNase is then added to digest away non-hybridized probe RNA. Labeled monoclonal anti-hybrid antibody is added to the tube. After incubation, excess RNase and conjugate is discarded and the hybrids detected.\nThe Digene test has some advantages: simplicity, rapidity, and quantitation of HPV. One major drawback of the Digene test appears to be difficulty of identification of the HPV type since mixtures of RNA probes for different HPV types are used. This is significant in view of the fact that there are many types of cancer-causing HPV types and current HPV vaccine will not protect humans from all types of HPV that cause cervical cancer. Therefore, the test is inadequate for complete identification of HPV type (s) in the infected cervical cells.\nOn Jun. 8, 2006, the FDA approved GARDASIL®, a prophylactic HPV vaccine which is marketed by Merck. The vaccine shows protection against initial infection with HPV types 16 and 18. Since the current vaccine will not protect women against all the HPV types that cause cervical cancer, it will be important for women to continue to seek Pap smear testing and other types of testing even after receiving the vaccine.\nVarious methods are available for identification of HPV type (genotyping). For example, a labeled HPV type-specific probe (labeled with radioisotope, fluorescence, biotin, etc) is directly hybridized with cellular DNA (usually after extraction and purification) in several formats. Liquid hybridization (such as the Digene “hybrid capture”), Southern and dot blot hybridization, and fluorescent in situ hybridization (FISH) are non-limiting examples. Gene amplification methods are also used. For example, HPV DNA is amplified by PCR by use of type-specific primers. However, this method suffers from the difficulty of identification of different types of HPVs in one experiment. The amplified DNA (usually labeled during PCR) can be used for hybridization with type-specific capture probes in dot blot, microtiter plate hybridization, or line probe assay in a similar way as DNA chip. However, these methods lack sensitivity and suffer from difficulties in data interpretation. DNA chip offers possibility of identification of many different types of HPV at once.\nRecently, Biomedlab (Seoul, Korea) has offered a DNA chip that has capture probes for as many as 30 different types of HPVs on a single chip that can distinguish low-risk as well as high-risk HPVs with high specificity and sensitivity (U.S. Pat. No. 7,301,015). The use of this chip requires purification of DNA from patient cervical cells, PCR and labeling of amplified HPV DNA, and hybridization of the amplified DNA with the type-specific capture probes immobilized on glass slide surface.\nCitation of the documents herein is not intended as an admission that any is pertinent prior art. All statements as to the date or representation as to the contents of these documents is based on the information available to the applicant and does not constitute any admission as to the correctness of the dates or contents of these documents."}
{"text": "1. Field of the Invention\nThis invention relates generally to rolling mills in which bars of a predetermined length are gathered into bundles, and is concerned in particular with a system and method for separating and recovering unacceptably short bars in advance of the bundling operation.\n2. Description of the Prior Art\nDuring the process of rolling bars in a bar mill, merchant mill, or section mill, short bars are invariably produced in the course of producing full length bars (known as commercial order lengths). The bar products can be ribbed reinforcing bars, plain rounds, squares, hexagonals, flats, angles, channels or even in some cases special shapes. Mill operators employ various methods to separate the short bars from the full length bars. These include optimization of billet length, optimization of rolled product length, cropping/chopping off of the short length arisings at mill shears or at the cold shear after the cooling bed. Each of these methods results in some form of cost added to end product.\nThe principal objective of the present invention is to provide a system and method for minimizing such added costs by separating and recovering short bars in a manner that improves upon the yield of the mill without adversely affecting the mill's throughput."}
{"text": "1. Field of the Invention\nThis invention relates generally to games, and more particularly to ball games.\n2. Description of the Prior Art\nMany ball games are not well suited for play in open areas. The balls roll far from the players and must be continually retrieved. This greatly detracts from the fun of the game and is also dangerous. Children chasing the ball sometimes run into streets where passing automobiles threaten their safety.\nSeveral games have been devised in which players throw a ball at a target. These games include basketball. Basketball requires a net mounted ten feet above the playing surface, however, and the ball must be thrown with some force to reach this height. The basket and backboard should be rigidly mounted to a heavy support, and an adequate basket and rim for basketball is generally not portable. Some players also do not have the physical ability to comfortably throw the ball to this height, and many young children, women, and older individuals cannot enjoy playing the game.\nPlaying areas sometimes are not paved, or are uneven and rough. This is especially the case in playgrounds, yards, and campsites. Games which require players to run are dangerous to play on these surfaces. Some games require players to bounce a ball on the playing surface, and these games are difficult or impossible to play on rough or uneven surfaces.\nIt would be desirable to provide a ball game which is suited for play in open areas. The game should preferably be designed to prevent children from chasing balls into streets. It further would be desirable if the game were fully portable and playable on uneven and rough surfaces. The game could thereby be played in playgrounds, yards, and campsites. The game preferably could be enjoyed by players with widely varying physical capabilities."}
{"text": "This invention is related to hybrid circuit chips or silicon substrate circuit flip chips which are soldered to a support structure. Specifically, there is provided a heating means integral with the circuit chip to permit unsoldering a specific chip which is soldered to a support structure to permit soldering a chip to a support structure without adversely affecting a neighboring chip and/or printed electrical conductors on the support structure.\nHybrid circuits usually include a plurality of circuit chips having one or more electrical components. These circuit chips are usually provided with a plurality of bonding pads for soldering to corresponding bonding pads on a glass or ceramic substrate which provide a support structure for the circuit chips. The hybrid circuit substrate or support structure also has a plurality of electrical conductors printed thereon using, typically, thin film or thick film techniques. The electrical conductors provide electrical interconnections between the circuit chips which have been soldered to the support structure.\nMany circuit chips are silicon monolithic integrated circuits. This variety of circuit chips has bonding pads on one surface thereof and includes \"solder bumps\" on the pads. This type of chip is sometimes referred to as a flip chip. The solder bumps are part of the soldering process for soldering the circuit chip to the hybrid circuit support structure as is well known.\nSince there is a very high cost in a completely assembled hybrid circuit, it is desirable to be able to repair an assembled hybrid circuit by removing a defective circuit chip and replacing it with an operative unit rather than scrapping the whole hybrid circuit. Prior art techniques for removing circuit chips from hybrid circuits include using a hand soldering iron, and using a rework fixture employing a hydrogen flame as a heating mechanism. These techniques, in many situations, result in irreparable damage to the electrical conductor patterns on the hybrid circuits as well as to adjacent good circuit chips and components. Another technique for removing hybrid circuits is shown in U.S. Pat. No. 3,904,100 in which a fixture is used in repairing hybrid circuits. The fixture includes a mechanism for heating a specific circuit chip to melt the soldering material between the circuit chip and the hybrid circuit substrate and permit removal of the circuit chip without damaging either adjacent circuit components or the conductor pattern on the hybrid circuit substrate. The latter technique for removing hybrid circuits has a disadvantage insofar as being a cumbersome apparatus to operate and requires an undue expense for providing the intended function."}
{"text": "This invention relates to a process for making expandable styrenic polymers anti-lumping during preexpansion by coating phosphoric acid esters of polyethoxyethanol ethers onto the expandable polymer particles by a dry blending process.\nThe making of low density, cellular, shaped, plastic articles from expandable particles of styrene polymers is well known. Such particles generally contain a blowing agent which boils below the softening point of the polymer and which will cause the particles to expand when they are heated. When the expanded particles are heated in a mold cavity, the particles expand further to fill the mold and fuse together to form a shaped article.\nThe formation of molded articles from expandable styrene polymer particles is effected in two steps: (1) pre-expanding of foamable styrene polymer particles to a density of between 1 and 3 pounds per cubic foot; and (2) further heating the pre-expanded particles in a closed mold to cause further expansion and form a fused, one piece article having the shape of the mold. The second step is what is normally referred to as \"molding\".\nThe pre-expansion step may be carried out by heating the expandable polymer particles by any suitable heating medium such as steam, hot air, hot water, or radiant heat. A widely used method of pre-expansion of the particles is a process such as that disclosed in U.S. Pat. No. 3,023,175 by Rodman.\nAn undesirable result of the pre-expansion before introduction into the mold is the tendency of the pre-expanded particles to clump together and form lumps which render the particles unsuitable for molding. These lumps cannot be properly conveyed in processing plants and are unsuitable for charging into molds of intricate contour due to improper fill-out of the molded articles. Even if reasonable fill-out is achieved, variations in density within the moled articles can result and voids can occur.\nAnother problem is that of stratification of prepuff particles in the continuous pre-expander which clogs the opening and necessitats periodic shut-down of the pre-expander to allow clean out of the stratified, densified chunks of foam.\nVarious methods have been proposed for the prevention of lumping and stratification during pre-expansion, but these have all suffered from some objectionable feature such as rapid loss of blowing agent, poor flow of the treated beads, dust problems which cause plugging of the steam ports, and often serious reduction in fusion of the particles on molding. U.S. Pat. No. 3,520,833 teaches the addition of lecithin during the impregnation of the particles with the blowing agent. Unfortunately, the lecithin imparts an undesirable odor to the molded articles. U.S. Pat. No. 3,462,293 teaches to coat the particles with polymeric materials by a fluid bed process. This process involves an additional expense of fluidizing the particles and coating with the polymer latexes. U.S. Pat. No. 3,444,104 teaches the addition of calcium silico aluminate. This additive tends to foul transfer lines and molds with clay-like deposits."}
{"text": "This invention relates to a process for the stereospecific polymerization of .alpha.-olefins and a catalyst therefor.\nMethods of supporting the titanium compound of Ziegler-Natta catalysts on a carrier to improve the catalyst activity have been developed in a commercial scale for polymerization of .alpha.-olefins and generally, have being employed for a catalyst for polymerization of ethylene.\nHowever, in case of polymerization of .alpha.-olefins such as propylene and butene, a useful, crystalline polymer can not be obtained unless the polymeric chain has an isotactic structure in which alkyl groups such as methyl or ethyl are stereospecifically regulated. Accordingly, catalysts improved only in polymerization activity as the case of ethylene polymerization cannot be considered to be a useful catalyst for polymerization of .alpha.-olefins and thus, regulating the stereospecificity of polymer is a significant matter.\nIn connection with this, there is, for example in British Pat. No. 1,435,768 provided a catalyst comprising a titanium composition, an organoaluminum compound and an electron donative compound, said titanium composition being prepared by copulverizing magnesium halide, a solid organic substance, an organic acid ester and titanium halide. Also, Japanese Patent Publication No. 52-50037 discloses a catalyst comprising a titanium composition and an organoaluminum compound, said titanium composition being obtained by reacting a copulverized product of magnesium halide and an organic acid ester with titanium tetrachloride. These catalysts are, however, still insufficient in a polymerization activity and a crystallinity of polymers obtained."}
{"text": "This invention relates to a process for the production of alkali metal silicate organic plastics by emulsifying a polymerable unsaturated organic compound and an organic epoxide compound with an aqueous alkali metal solution by mixing the polymerable organic compound and an organic epoxide compound with an aqueous solution of alkali metal silicate then adding a salt-forming compound in the amount up to 10%, based on the alkali metal silicate, preferably an organic acid, while agitating thereby producing a stable emulsion. A polymerizing catalyst such as a peroxide type catalyst is added to the emulsion thereby producing a poly(alkali metal silicate-polymerable organic compound epoxide compound) copolymer. In most products an excess amount of the aqueous alkali metal silicate may be used. The inorganic-organic plastic produced by the process of this invention has greatly improved flame resistance properties.\nThe polymerization of an alkali metal silicate with a polymerable unsaturated organic compound was illustrated in U.S. patent application Ser. No. 71,628, filed Sept. 11, 1970, by David H. Blount. The alkali metal silicate is oxidized by a peroxide initiator then polymerized with a polymerable organic compound. I have discovered that a stable emulsion of an aqueous alkali metal silicate and a polymerable unsaturated organic compound may be produced by adding up to 10% by weight, percentage based on weight of the aqueous alkali metal silicate solution, of a salt forming compound, and mixing with the mixture of the aqueous alkali metal silicate and polymerable compound. This stable emulsion greatly enhances the reaction between the alkali metal silicate and polymerable organic compound. Any suitable polymerable unsaturated organic compound may be used in this invention that can be polymerized in an aqueous alkali metal silicate solution in the presence of a peroxide initiator.\nThe emulsions of inorganic-organic plastics may be used as an adhesive on wood, paper, cement, plastics, ceramics, etc., as a coating agent on wood, cement, plastics, ceramics, etc., and may be dried or coagulated with a salt forming compound to produce a molding powder which may be molded by heat and pressure to produce useful objects such as knobs, handles, gears, pipes, toys, etc. The emulsion of inorganic-organic plastics may be further reacted with organic compounds such as polyisocyanates, isocyanates, epoxide compounds, substituted organic compounds, water-binding agents and many other compounds. The emulsion of inorganic-organic plastics may be used as a cavity filler, as putty, as a caulking compound, and in producing laminates.\nIt is accordingly, an object of my invention to provide novel inorganic-organic copolymers. A further object is to provide novel copolymers which may be used as an adhesive. A further object is to provide novel copolymers that will react with polyisocyanates to produce useful resinous and foam products. A further object is to provide a process for preparing novel inorganic-organic copolymers. Another object is to produce emulsions of inorganic-organic copolymers which may be used to produce concrete reinforced and reacted with inorganic-organic copolymers.\nThe inorganic-organic plastics may be produced by emulsifying and reacting the following components:\nComponent (a) an aqueous alkali metal silicate solution; PA0 Component (b) a polymerable unsaturated organic compound; PA0 Component (c) a salt forming compound; PA0 Component (d) an initiator PA0 Component (e) an epoxide compound PA0 p,p'-diphenylmethane diisocyanate PA0 phenylene diisocyanate PA0 chlorophenylene diisocyanate PA0 tolylene diisocyanate PA0 m-xylylene diisocyanate PA0 benzidine diisocyanate PA0 naphthylene diisocyanate PA0 tetramethylene diisocyanate PA0 pentamethylene diisocyanate PA0 hexamethylene diisocyanate PA0 decamethylene diisocyanate PA0 thiodipropyl diisocyanate"}
{"text": "The present invention relates to storm windows and, more particularly, to an improved storm window frame unit having an angulated, peripheral flange for direct interengagement with an existing window jamb.\nEnergy conservation in residential housing is of prime import in contemporary society. In that respect, conventional single glaze windows have been shown to be the biggest single cause of energy loss in the typical home. Energy is lost first through air infiltration around the window frame and sash. Window energy loss is also predominant via radiation through the glass in the window lites. It has been shown that approximately three (3) times as much energy is lost via radiation through the window lite as is normally lost from air infiltration. This is particularly so because one square foot of glass conducts at least ten times as much heat as one square foot of conventional wall. Glass has an R factor, or thermal resistance value, of 0.88. A conventional insulated wall provides a resistance, or R factor of 13 or better. Overall conventional windows have been shown to exhibit problems in over 80 percent of all homes. This figure exceeds half a billion windows without including the window problems of office buildings and other types of commercial establishments. Such an energy problem accounts for a great fuel loss which under the present fuel crisis renders the correction of such problem matters of national socio-economic import.\nConventional window construction include both single and double hung windows, the frames of which are more conventionally constructed of steel or aluminum of the extruded variety. Most conventional homes include windows of single glass panes which, although usually effectively sealed therearound readily permit the radiation of energy therethrough. Such heat loss through radiation leads to a great energy loss and an aggravation of the energy shortage. It is thus common in present times to encounter improved window designs including double pane window constructions utilizing insulated thermal barriers between glass panes for reducing conduction heat loss. The only effective means for reducing radiation heat loss is the utilization of blinds, curtains and the like for covering the transparency of the associated window.\nAnother form of reducing conduction and conduction heat loss in existing single pane windows is the utilization of storm windows being affixed at the outward regions thereof for creating a dead air space therebetween. According to the National Bureau of Standards, storm windows properly installed will cut heating bills a minimum of 10 to 20 percent. At average 1977 fuel prices in a 6000 degree day climate, such as Cleveland and Chicago, the investment in storm windows which are properly installed, can be recovered in five years. Unfortunately, only about 23 million of the 75 million yearround, single family homes in the United States have storm windows on every window, as reported by the Census Bureau. Such storm windows are commonly affixed directly to the frame or jamb of conventional windows for creating a sealed dead air space therebetween. Most prior art construction includes an insulation member such as box vinyl disposed between the existing window jamb and the window frame of the storm window. The insulation member reduces heat transfer and serves as a spacer for separating the two windows.\nConventional storm windows are fabricated from extruded aluminum facilitating the fabrication of multiple slide channels for receiving a plurality of sashes. The storm window is affixed to the existing window jamb by the utilization of screws or the like extending through an outer flange formed around the storm window frame. The flange is constructed in generally parallel relationship with the integrated window channel units of the storm window frame which peripheral extension in parallelism necessitates the use of a channel insulating element between the storm window and the conventional window for separation therebetween. The separation facilitates the creation of a dead air space as well as eliminating interference between the moving parts of the conventional window and the storm window. Without the separation the use and effectiveness of the window would be lessened.\nCertain problems inherent in the fabrication and installation of conventional storm windows is the outer securement flange which usually necessitates the incorporation of an insulated channel member thereunder when applied to conventional windows. This channel member must be extruded from plastic or similar insulative material and include a flange configuration in itself for sealably engaging the subject storm window flange therewithin. The channel must also be of sufficient rigidity to withstand the compression and intrusion of fasting members therethrough. In short, most prior art storm window construction necessitates the creation of a suitable insulative channel and the disadvantages of securing the window thereupon which provides numerous areas of installation error and insulation ineffectiveness. Such problems may cause air infiltration and related energy loss of the type sought to be eliminated by the use of the storm window in the first place. It would be an advantage, therefore, to provide a storm window in a frame having a configuration which would eliminate the need for insulative channels and the related problems associated therewith.\nThe storm window of the present invention is provided for just such an application wherein the extruded frame of the storm window includes the angulated flange portion extending away from the extruded channel portion for purposes of installation and operation. The angulated flange includes a generally coplanar flange therearound for direct abutting engagement with the window jamb. In this manner, a substantially flat insulating gasket or plasticized sealant may be utilized between the existing window jamb and the storm window flange instead of a separate type insulating member of the prior art box vinyl variety. In this manner, the installation of the storm window of the present invention is greatly facilitated and the effectiveness of the sealing therebetween is maximized."}
{"text": "Conventional user interfaces for appliances include mechanical switches, pushbuttons, indicator lights, clock display modules, potentiometers and/or electronic controls. Generally, numerous device elements can be used to communicate numerical and non-numerical information.\nOne conventional user interface includes a seven segment display module for communicating any numerical value. For example in a range cooktop, the seven-segment display may display the clock value, a timer value and/or a temperature value. In a washing machine, dryer or dishwasher, the seven segment display may indicate the amount of time remaining in a cycle. In addition to a seven segment display, conventional user interface further include thru-hole light emitting diodes (LEDs) to communicate any other information that is not number dependent. For example, in a cooktop which burner is activated, still hot, or an indication of oven activation. In a washing machine, dryer or dishwasher, the selected cycle or temperature may be communicated by a thru-hole LED. This configuration of a user interface is undesirable because the seven segment display module has a large footprint that increases the size of the overall user interface and two the different elements increase the manufacturing time and costs.\nAnother conventional user interface replaces the seven segment display and through-hole LEDs with a liquid crystal display (LCD). All user information is communicated by the LCD. This configuration of a user interface is undesirable because it increases the programming requirements and draws more power. In addition, when this type of display becomes inoperable, the entire appliance may be inoperable.\nTherefore, a need exists for a system and method of displaying information to a user that overcomes the above mentioned disadvantages. A user interface that could be manufactured using the same elements for the number information as well as any other information communicated to the user would be useful. In addition, it would be advantageous to use reverse mount LEDs to improve design flexibility, reliability and reduce the size and cost of the user display."}
{"text": "Electronic devices such as mobile phones, laptop computers, electronic notebooks, digital cameras, and video game instruments usually have an audio function. Such an electronic device having an audio function includes a housing, in which is placed a sound emitting part such as a speaker or a buzzer or a sound receiving part such as a microphone. The housing is provided with an opening at a position corresponding to that of the sound emitting part or the sound receiving part. Sounds are transmitted through the opening. A sound-transmitting membrane is attached to the housing so as to cover the opening of the housing in order to prevent foreign matters such as water droplets from entering the housing. The sound-transmitting membrane allows sound to pass through but prevents foreign matters from passing through. Known examples of the sound-transmitting membrane include porous plastic membranes such as porous polytetrafluoroethylene membranes and porous ultrahigh molecular weight polyethylene membranes (see Patent Literature 1).\nPatent Literature 2 describes a waterproof sound-transmitting member including a sound-transmitting membrane and a housing-side adhesive layer. The housing-side adhesive layer is a layer used to bond the waterproof sound-transmitting member to the housing and is laminated on the sound-transmitting membrane.\nPatent Literature 3 describes a waterproof sound-transmitting member including a sound-transmitting membrane, a support layer, and a housing-side adhesive layer. In the waterproof sound-transmitting member described in Patent Literature 3, the entire peripheral portion of the support layer extends outwardly beyond the housing-side adhesive layer."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for producing an electrostatic image developing toner which is useful in used electrophotographic system copying machines and printers. More particularly, the invention relates to a method for producing an electrostatic image developing toner, which comprises conducting particle size control in an agglomeration step to prepare toner particles having a sharp particle size distribution with a lower content of fine powder and coarse grains.\n2. Discussion of Background\nA conventional electrostatic image developing toner which has been widely used in the field of electrophotography is produced by melt-kneading a mixture comprising a coloring agent such as carbon black or a pigment, wax and/or a magnetic material with a resin (polymer) by an extruder, followed by pulverization and classification. However, by using this method, it is difficult to disperse the coloring agent and the compounding agent in the resin.\nOn the other hand, in recent years there has been a need for printers and copying machines having improved performance with respect to high picture quality and a high printing speed. To achieve a high picture quality, the toner is required to consist of small particle sizes ranging from 3 to 9 xcexcm and a narrow particle size distribution. Further, to achieve a high printing speed, the toner is required to have a high fixation rate, i.e. it should be able to be fixed at a low temperature, as a toner physical property. Such are also favorable from the viewpoint of providing an environmental printer or copying machine which shortens the standby time, which does not generate much heat, and which has a small electric power consumption.\nIn order to achieve fixation of the toner at a low temperature, the softening point of the resin should be lowered. However, the conventional toner obtained by the melt-kneading/pulverization method has an antinomic relation such that when the softening point of the resin is lowered, the fixation property at low temperature will be improved, but at the same time, the Tg of the toner tends to be low, and the blocking resistance tends to be impaired. Namely, it is impossible to control the dispersion of the materials, and the structure control which satisfies both the blocking resistance and the fixation property at low temperature is extremely difficult.\nA method of mixing an additive such as wax in the toner formulation is known, but the added amount of the additive is limited to a level of from 4 to 5 parts by weight based on 100 parts by weight of the resin, and the fixation characteristic of the toner at a low temperatures tends to be inadequate when the additive is incorporated in the toner in such amounts. Further, when flakes obtained by melt-kneading are mechanically pulverized to prepare a toner, the yield tends to be poor, and the particle size distribution tends to be broad.\nParticularly when it is attempted to obtain a toner having small particle sizes, this tendency is very significant.\nIn recent years, the known production method of producing polymer toner by-means of emulsion polymerization or suspension polymerization is expected to take the place of the melt-kneading/pulverization method. By these methods, it is possible to control the dispersion of the materials in a manner which is different from the melt-kneading/pulverization method. Further, it is also possible to prepare a toner having small particle sizes and a good particle size distribution. Particularly by the emulsion polymerization agglomeration method, it is possible to control the particle size, the particle size distribution and the shape of the toner.\nIn the case of producing a toner by the emulsion polymerization method, a pigment and a charge controlling agent are added to a dispersion of an emulsified resin, containing primary resin particles having particle sizes of ranging from 0.05 to 0.5 xcexcm obtained by polymerization, and also e.g. an electrolyte to agglomerate the primary particles thereby preparing a toner of a particle size ranging from 3 to 9 xcexcm. The slurry obtained containing the toner particles is washed and dried thereby preparing toner particles as a product. However, it is difficult to prepare toner particles having a uniform particle size in the agglomeration step even by this method. Stirring plays an important role in the agglomeration of the primary particles. That is, in the agglomeration step, the viscosity of the mixture changes as the particles agglomerate. Accordingly, it is necessary to control the agglomeration of particles by changing the-rate of stirring the mixture as the viscosity of the medium changes during agglomeration. However, no method has been proposed which adequately controls the agglomeration of particles. Accordingly, it is difficult prepare a uniform toner having small particle sizes even by emulsion polymerization/agglomeration, and it is necessary to further e.g. process the agglomerated particles such as by classification after the agglomeration step to achieve the desired performance of the toner.\nAccordingly, one object of the present invention is to provide a method of preparing toner particles having small particle sizes and a sharp particle size distribution, suitable to achieve a high picture quality and a high printing speed, by agglomerating primary polymer particles prepared by e.g. emulsion polymerization.\nIt has now been found that the above objective can be achieved by controlling the stirring conditions during agglomeration of primary polymer particles under specific conditions with a view to uniformly mix the liquid containing the particles for agglomeration so that the agglomeration medium thickens during agglomeration, provides sharpness of the particle size distribution of the agglomerated particles and provides ease of particle size control.\nBriefly, this object and other objects of the present invention as hereinafter will become more readily apparent can be attained by a method for producing an electrostatic image developing toner, which comprises, stirring a dispersion containing at least primary polymer particles and coloring agent particles in a stirring tank, thereby agglomerating the particles, wherein the shape of the liquid level during the stirring operation in the agglomeration step satisfies the following formulae (I) and (II):\nHCxe2x89xa60.8Hxe2x80x83xe2x80x83(I), \n1.2Hxe2x89xa6HExe2x80x83xe2x80x83(II) \nH: the distance between the bottom of the stirring tank at its center and the horizontal plane of the surface of the dispersion before the start of stirring,\nHC: the distance between the bottom of the stirring tank at its center and the liquid level of the surface of the dispersion at its center during stirring, and\nHE: the distance between the bottom of the stirring tank at its center and the liquid level of the surface of the dispersion at the edge of the tank during stirring.\nAnother aspect of the present invention is a method for producing an electrostatic image developing toner, which comprises, an agglomeration step of stirring a dispersion containing at least primary polymer particles and coloring agent particles in a stirring tank to prepare agglomerates of the primary particles and coloring agent particles, wherein a double helical ribbon blade is employed as the agitating blade within the stirring tank."}
{"text": "1. Technical Field\nAn improved manifold/nozzle closure system for fluid dispensers is shown and described. The disclosed closure system is a manual system that, in a closed position, provides a cover/closure and a seal for a manifold/nozzle through which one or more fluids are dispensed. In the closed position, the closure element is disposed beneath the manifold or nozzle where it collects any fluid drippings between dispensing operations and provides a sealing effect to prevent dried material from clogging the nozzle(s). To move the closure system to an open or to a dispense position, an easily accessible push plate is provided which, when pushed away from the user and towards the machine housing, moves the closure element downward before pivoting the closure element laterally rearward and away from the nozzle or manifold. After the fluid is dispensed, a biasing element returns the closure element back to a sealing engagement with the manifold/nozzle. The closure element preferably includes a removable cup and a soft seal, either of which may be disposable and/or cleanable.\n2. Description of the Related Art\nSystems for dispensing a plurality of different fluids into a container have been known and used for many years. For example, systems for dispensing paint base materials and colorants into a paint container are known. These paint systems may use twenty or more different colorants to formulate a paint mixture. Each colorant is contained in a separate canister or package and may include its own dispensing pump. The colorants and the respective pumps may be disposed on a turntable or along one or more horizontal rows. In a turntable system, the turntable is rotated so that the colorant to be dispensed is moved to a position above the container being filled. In designs using one or more horizontal rows, the container may be moved laterally to the appropriate colorant/pump.\nSystems for dispensing large varieties of different fluids are not limited to paints, but also include systems for dispensing pharmaceutical products, hair dye formulas, cosmetics of all kinds, nail polish, etc. Some systems for use in preparing products at a point of sale may use a stationary manifold through which a plurality of nozzles extend. Each fluid to be dispensed is then pumped through its own individual nozzle that is accommodated in the manifold. Depending upon the size of the container and the quantity of the fluids to be dispensed, manifolds can be designed in a space efficient manner so that a single manifold can accommodate twenty or more different nozzles. The nozzles are connected to the various ingredients by flexible hoses and the ingredients are contained in stationary canisters or containers.\nIn many fluid dispensing applications, precision is essential as many formulations require the addition of precise amounts of certain ingredients. This is true in the pharmaceutical industry but also in the paint and cosmetic industries as the addition of more or less tints or colorants can result in a visible change in the color of the resulting product.\nOne way in which the precision of dispensing systems is compromised is “dripping.” Specifically, a “leftover” drip may be hanging from a nozzle that was intended to be added to a previous formulation and, with a new container in place under the nozzle, the drop of liquid intended for a previous formulation may be erroneously added to a new formulation. Thus, the previous container may not receive the desired amount of the liquid ingredient and the next container may receive too much.\nTo solve the drip problem, various scraper and wiper designs have been proposed. However, these designs often require one or more different motors to operate the wiper element and are limited to use on dispensing systems where the nozzles are separated or not bundled together in a manifold. Use of a wiper or scraping function would not be practical in a multiple nozzle manifold design as the ingredients from the different nozzles will be co-mingled by the wiper or scraper which would then also contribute to the lack of precision of subsequently produced formulations.\nAnother problem associated with dispensing systems that make use of nozzles lies in the dispensing of relatively viscous liquids such as tints, colorants, base materials for cosmetic products, certain pharmaceutical ingredients or other fluid materials having relatively high viscosities. Specifically, the viscous fluids have a tendency to dry and cake onto the end of the nozzles, thereby requiring frequent cleaning in order for the nozzles to operate effectively. While some mechanical wiping or scrapping devices are available, these devices are not practical for multiple nozzle manifold systems and the scraper or wiper element must be manually cleaned anyway.\nOne solution provided in commonly assigned US 2005/0205154 is a mechanized closure element that covets and seals from beneath the manifold after the dispensing operation is complete. In this manner, the viscous materials being dispensed through the nozzles have less exposure to air thereby requiring a lower frequency of cleaning operations However, the mechanized system of US 2005/0205154 requires modification of the existing system software. Thus, it is relatively expensive and is difficult to add to an existing system as a retrofit. Thus, a more economical approach is needed."}
{"text": "The present invention relates to a fastening device for a line, in particular for an electric cable in a motor vehicle.\nDE 43 41 384 C2 describes a feedthrough for lines, which includes elastic grommets which are lockable in place in a carrier and accommodate the line in a sealed manner. The elastic grommets may be slid via their base regions in a fitted manner into recesses which are connected to the edge of the carrier. The recesses for the grommets are radial in design and have a width which corresponds to the diameter of the grommet. A ring which encloses the carrier on the edge is also provided in order to lock the grommets in place in the recesses. The ring secures the grommets in the inserted position. The entire system has a complex design, and inserting the cables with the grommets is time-consuming."}
{"text": "The fabrication of modern integrated circuits (ICs) involves several process steps. Integrated circuits are first fabricated on/in a semiconductor wafer. The semiconductor wafer contains multiple duplicated semiconductor chips, each including integrated circuits. The semiconductor chips are then sawed from the semiconductor wafer and packaged. The packaging process has two main purposes: to protect delicate semiconductor chips, and to connect interior integrated circuits to exterior connections.\nAs semiconductor devices reach higher levels of integration, the packaging of a semiconductor chip accounts for a considerable portion of the cost of producing the device. Failure of the package leads to costly yield reduction.\nThere are challenges in fabricating reliable package structures."}
{"text": "This application is the national phase under 35 U.S.C. xc2xa7371 of PCT International Application No.\nPCT/JP99/00391 which has an International filing date of Jan. 29, 1999, which designated the United States of America.\nThe present invention relates to a method for manufacturing a pseudo-boehmite powder as a raw material for producing a carrier for a catalyst for hydrogenation refining. In particular, the present invention relates to a method for manufacturing pseudo-boehmite by effectively utilizing waste water and sludge discharged from the aluminum surface-processing steps, and the present invention also relates to a method for producing a catalyst carrier and a method for producing a hydrogenation catalyst.\nIn the petroleum refining, a catalyst, in which a metal component having a hydrogenation ability is carried on an alumina carrier, is used as catalyst for hydrogenation refining in many cases. The alumina carrier can be produced by mixing, kneading, and forming a fine powder of pseudo-boehmite as an alumina hydrate, followed by calcination. The pseudo-boehmite powder, which is used as a raw material for the alumina carrier as described above, is produced by using a first method (1) in which an aluminic acid or aluminum salt is formed from gibbsite (aluminum hydroxide) followed by performing a reaction therewith, or a second method (2) in which aluminum alkoxide is hydrolyzed to synthesize the pseudo-boehmite powder. The characteristics of the obtained catalyst are greatly affected by the quality and the characteristics of the alumina carrier and the pseudo-boehmite. Therefore, in general, a high purity aluminum source is used as the raw material of the pseudo-boehmite powder.\nWhen an aluminum product to be used for building materials such as sashes made of aluminum is produced, a treatment is performed in which an aluminum oxide coating forms on the formed aluminum by means of an electrochemical method to improve the corrosion resistance, the abrasion resistance, and the appearance of beauty. This treatment uses a process comprising the steps of, for example, anodization, etching, and washing with water, and it is called the alumite treatment. FIG. 1 shows a series of aluminum surface treatment process including the alumite treatment.\nAs shown in FIG. 1, the aluminum material is first subjected to the surface degreasing, and it is washed with water. After that, the material is subjected to the alkaline etching step. In the alkaline etching step, the surface of aluminum is dissolved with caustic soda or sodium hydroxide in an amount of several xcexcm to several tens xcexcm. Subsequently, the steps of washing with water, neutralization, and washing with water are performed, and then the material is subjected to anodization or an anodic oxidation step. In the anodization, an oxide coating is formed on the surface of the aluminum material in an electrolytic solution such as sulfuric acid. After the anodization, the washing step is further performed with water, optionally followed by coloring and sealing steps to obtain the product.\nA large amount of alkaline waste liquid, in which metal aluminum is dissolved at a dense concentration, is produced in the alkaline etching step. In order to effectively utilize the alkaline waste liquid, aluminum hydroxide is deposited from the alkaline waste liquid. On the other hand, the resultant recovery liquid is reused by returning it as caustic soda to the alkali etching step. The deposited aluminum hydroxide has its water concentration of 12 to 15%, and it is a crystalline aluminum hydroxide called gibbsite. This product has been hitherto used as a material for cement and ceramic industries.\nIn the anodization step described above, an acidic waste liquid (sulfate waste liquid) dissolved with metal aluminum is also produced. In order to treat the acidic waste liquid, diffusion dialysis method has been hitherto used so that sulfuric acid and aluminum sulfate are recovered. In other cases, an alkaline waste liquid, which is produced by the alkaline etching, has been hitherto mixed with the acidic waste liquid to recover aluminum hydroxide.\nFurther, aluminum dissolved in a washing liquid also exists in the step of washing with water after the alkaline etching and the steps of washing with water before and after the anodization. The washing liquids are gathered and neutralized, followed by solid-liquid separation to obtain sludge or slime containing aluminum hydroxide. The aluminum hydroxide in the sludge has a high water content, and it is an amorphous aluminum hydroxide in a form of gel. Therefore, the aluminum hydroxide sludge has troubles concerning the treatment and disposal. The sludge contains a large amount of other metal components which are components of the aluminum alloy other than aluminum, and it has been hitherto processed by reclamation, or it has been hitherto processed by dissolving it in a by-product of aluminum sulfate. The aluminum sulfate contains colored ions such as those of copper and nickel having been contained in the aluminum alloy. Therefore, the aluminum sulfate is inferior in added value as compared with the high purity aluminum sulfate produced by dissolving gibbsite. As a result, for example, a problem arises in that the aluminum sulfate can be merely used for usage as a neutralizing aggregation treatment agent for processing the waste water.\nThe effective utilization and the problems, which relate to the waste liquid and the aluminum hydroxide sludge in the alumite processing as described above, are disclosed in detail in xe2x80x9cTechnology for Processing Waste (Aluminum Hydroxide Sludge) and Reproducing Resource from Wastexe2x80x9d issued by the juridical foundation of Clean Japan Center (February, 1981). Factories for processing and treating alumite are present at not less than 1000 sites only in Japan. The amount of produced sludge exceeds 5000 tons per year. Therefore, the waste liquid, which is produced by the alumite processing or treatment, contains a large amount of valuable components. However, in the present circumstances, it is difficult to consider that such valuable components are sufficiently reused. Further, it is demanded to reuse aluminum scrap and aluminum waste produced from the alumite treatment and from other aluminum-processing sites.\nAn object of the present invention is to provide a method for reproducing a resource from waste, in which a waste liquid or a sludge containing aluminum as described above is utilized as a raw material for producing a catalyst carrier.\nAnother object of the present invention is to provide a novel process for producing pseudo-boehmite which serves as a raw material for a catalyst carrier. Still another object of the present invention is to provide a method for producing the catalyst carrier and a method for producing a hydrogenation refining catalyst based on the use of the novel process.\nStill another object of the present invention is to provide a novel method for treating alumite including the process of effectively utilizing a waste liquid and an aluminum hydroxide sludge produced from the alumite treatment.\nAccording to a first aspect of the present invention, there is provided a method for producing pseudo-boehmite, comprising the steps of dissolving metal aluminum in an acidic or alkaline aqueous solution to prepare an acidic or alkaline aluminum aqueous solution, and producing the pseudo-boehmite from the acidic or alkaline aqueous solution.\nAs described in the column of the background art, the pseudo-boehmite, which is the material for producing the alumina carrier for the catalyst, is generally produced by forming aluminic acid by using high purity gibbsite as a starting material, and then reacting the obtained product. The present inventors have found out that the pseudo-boehmite can be produced by reacting, under a predetermined condition, the alkaline aqueous solution or the acidic aqueous solution dissolved with a large amount of metal aluminum discharged in a large amount as a waste liquid from the aluminum-processing site such as those for the alumite treatment. The use of this method makes it possible to reuse aluminum as a resource which would be otherwise dealt with as a waste. Further, the use of this method makes it possible to produce, at low cost, the pseudo-boehmite as well as the catalyst carrier and the catalyst based on the use of the same. The method has not been hitherto known, in which the pseudo-boehmite is produced from the starting material of the acidic or alkaline aqueous solution of aluminum obtained by dissolving metal aluminum (metallic aluminum) as described above. In this specification, the phrase xe2x80x9cmetal aluminumxe2x80x9d of the phrase xe2x80x9cdissolve metal aluminumxe2x80x9d is intended to include oxide or hydrate formed by oxidation of metal aluminum is dissolved.\nIn order to form the pseudo-boehmite, an acidic aluminum aqueous solution may be mixed with an aluminate aqueous solution obtained from the alkaline aluminum aqueous solution described above. Alternatively, an aluminate aqueous solution may be mixed with the acidic aluminum aqueous solution described above. Those usable as the aluminate aqueous solution include, for example, an aluminate aqueous solution produced from the alkaline etching step for aluminum in the alumite treatment. Those usable as the acidic aluminum aqueous solution include an acidic waste liquid produced from the anodization step for aluminum in the alumite treatment, i.e., an acidic aluminum aqueous solution. The acidic aluminum aqueous solution may contain metal components other than aluminum in a molar ratio of 0.001 to 0.02 with respect to total aluminum contained in the aqueous solution.\nThose intended as the metal aluminum may include an aluminum material to be used in the aluminum-processing site such as those for the alumite treatment, because of the following reason. That is, a part of such an aluminum material is dissolved in alkali or acid in the following step such as the etching step to produce the acidic or alkaline aluminum aqueous solution which may serve as the material for producing the pseudo-boehmite.\nThe material may include an aluminum material used for any usage other than the alumite treatment.\nAccording to a second aspect of the present invention, there is provided a method for producing pseudo-boehmite, comprising the steps of dissolving metal aluminum in an alkaline aqueous solution to prepare an aluminate aqueous solution, dissolving metal aluminum in an acidic aqueous solution to prepare an acidic aluminum aqueous solution, and mixing the prepared aluminate aqueous solution with the prepared acidic aluminum aqueous solution to form the pseudo-boehmite. According to this method, for example, the pseudo-boehmite can be produced only by mixing, under a predetermined condition, the aluminate aqueous solution produced from the alkaline etching step for aluminum in the alumite treatment and the acidic aluminum aqueous solution produced from the anodization step for aluminum in the alumite treatment. Therefore, it is possible to establish an extremely low cost process for producing pseudo-boehmite by utilizing the waste liquid. This process can be carried out in the alumite treatment factory. Accordingly, this process provides the by-product for the alumite treatment factory.\nAccording to a third aspect of the present invention, there is provided a method for producing pseudo-boehmite, comprising the step of maturing or aging, in an alkaline solution, an aluminum hydroxide sludge produced from an alumite treatment. The present inventors have succeeded in producing the pseudo-boehmite only by treating the aluminum hydroxide sludge under a predetermined condition, while being limited to the waste liquid produced from the alumite treatment. According to this method, it is possible to effectively utilize a huge amount of aluminum hydroxide sludge discharged from the aluminum-processing sites existing in Japan as well as all over the world. Especially, when an alkaline waste liquid, which is produced from the alkaline etching step in the alumite treatment, is utilized as the alkaline solution described above, it is possible to produce the pseudo-boehmite in a closed system established at the alumite-processing site. The aluminum hydroxide sludge may be matured, for example, in the alkaline solution at a temperature of 50 to 90xc2x0 C. under a condition of pH 8 to 12. In general, the aluminum hydroxide sludge is produced from the alkaline etching step for aluminum, and the washing steps performed before and after the anodization step in the alumite treatment.\nAccording to a fourth aspect of the present invention, there is provided a method for treating alumite comprising the steps of alkaline etching for aluminum, anodization, and washing performed thereafter, the method further comprising the step of producing pseudo-boehmite by using an alkaline waste liquid produced from the alkaline etching step, an acidic waste liquid produced from the anodization step, or an aluminum hydroxide sludge produced from the washing step. In the alumite treatment method according to the present invention, the pseudo-boehmite, which serves as a raw material for a catalyst carrier, can be provided as a by-product. Especially, the pseudo-boehmite can be produced in the alumite treatment factory by mixing the alkaline waste liquid produced from the alkaline etching step and the acidic waste liquid produced from the anodization step. In addition to this process, or separately from this process, the pseudo-boehmite may be produced by maturing, in an alkaline solution, the aluminum hydroxide sludge produced from the washing step. The added value may be provided for the alumite treatment factory only by installing, at the alumite-processing site, the equipment for carrying out the process for producing the pseudo-boehmite as described above."}
{"text": "Data visualization applications enable a user to understand a data set visually, including distribution, trends, outliers, and other factors that are important to making business decisions. Some data sets are very large or complex, so the process of analyzing a data set, loading the data set, and displaying a corresponding data visualization can be slow. The process is also slow when a user chooses to change what data is displayed or how the data is displayed.\nData visualizations are often shared with others, sometimes in combination with other data visualizations as part of a dashboard. In some cases, the distributed data visualizations are static. To the extent a distributed data visualization or dashboard is dynamic, updates may be slow, particularly within a browser or on a mobile device."}
{"text": "Immune dysregulation can be a component of many pathological diseases or conditions. Such dysregulation may be a factor that favors the establishment, maintenance or progression of these diseases or conditions. Deficient immune responses or immune suppression can enhance a mammal's susceptibility to infection or to the development of cancer. Conversely, excessive or inappropriate immune responses can play a role in the establishment or progression of unwanted inflammation or autoimmune conditions. It would thus be advantageous to utilize agents that can modulate immune responses and to at least partially reverse immune dysregulation when such dysregulation is a component of a given pathological condition. Some agents are known that can ameliorate some aspects of immune dysregulation, but typically such agents have their own unwanted effects on either the host's immune system or other organs or tissues. Agents such as glucocorticoid steroids have been used to reduce unwanted inflammation in a number of clinical conditions, but such compounds often have serious limitations, such as inducing immune suppression or causing unwanted mood changes.\nHuman and mammalian immune responses to infections or other conditions are often characterized by responses mediated by different immune effector cell populations. In some situations, helper T cells designated Th1 in the murine system, facilitate immune effector functions that are typically dominated by cell-mediated responses. In other cases, helper T cells designated Th2 cells facilitate immune effector functions that are typically dominated by humoral responses. A vigorous Th1 response is usually desirable to help clear infections or to slow the progression of an infection. When a subject's immune response is biased to, or dominated by, a Th2-type response, the cytokines associated with the Th2 response tend to suppress the immune system's capacity to mount a vigorous Th1 response at the same time. The converse is also generally true. When mammalian immune responses begin to result in an increasing Th1 response, Th2 responses tend to weaken. Insufficient Th1 responses may be associated with progression of some infections or other conditions, see, e.g., M. Clerici and G. M. Shearer, Immunol. Today 14:107-111, 1993; M. Clerici and G. M. Shearer, Immunol. Today 15:575-581, 1994.\nMethods to make and use 3β-hydroxyandrost-5-ene-17-one (dehydroepiandrosterone) and certain other steroids and their biological properties have been described, see, e.g., U.S. Pat. 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Nos. 5,763,433 and 6,372,732 and PCT publication WO 96/35428 describe the capacity of certain androstane and androstene compounds such as 3β-hydroxyandrost-5-ene-17-one to treat certain immune disorder conditions such as systemic lupus erythematosus.\nU.S. Pat. Nos. 5,925,630, 5,939,545 and 5,962,443 describe the capacity of 19-nur-pregnane steroids, 3α-hydroxy-5α-pregnan-20-one and related steroids to modulate certain neurological activities such as hypothalamic function and GABA receptor activity.\nOther biological effects and/or metabolic conversions of steroid compounds have been described, e.g., Batta et al., J. Biol. Chem. 1986 25:127-133, Belli et al., Liver 1991 11:162-169, Bhattacharjee et al., Anal. Biochem. 1992 201:233-236, Blake et al., Int. J. Peptide Protein Res. 1982 20:97-101, 1986 25:127-133, Bonaventura, Am. J. Obstet. Gynecol. 1978 131:403-409, Bucala et al., J. 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Lipid Res. 1995 36:628-638, Jellinck et al., Steroids 1967 10:329-346, Jonsson et al., J. Pediatr. Gastroenterol. Nutr. 1995 20:394-402, Kalimi et al, Mol. Cell. Biochem. 1994 131:99-108, Kramer et al., J. Biol. Chem. 1994 269:10621-10627, LaRochelle et al., Steroids 1984 43: 209-217, Liao et al., Carcinogenesis 1998 19:2173-2180, Lillienau et al., J. Clin. Invest. 1992 89:420-431, Loria, Psychoneuroendocrinology 1997 22:S103-S108, Luscher et al Mol. Immunol. 1983 20:1099-1105, Manna et al., J. Biol. Chem. 1999 274:5909-5918, Marschall et al., J. Biol. Chem. 1989 264:12989-12993, Medh et al., Cancer Res. 1998 15:3684-3693, Mohan et al., Steroids 1992 57:244-247, Munoz de Toro et al., J. Steroid Biochem. Mol. Biol. 1998 67:333-339, Padgett et al., J. Neuroimmunol. 1998 84:61, Padgett et al., Ann. N.Y. Acad. Sci. 1995 774:323, Padgett et al., J. Immunol. 1994 153:1544-1552, Pashko et al., Carcinogenesis 1984 5:463-466, Pashko et al., Carcinogenesis 1981 2:717, Petrylak et al., J. Clin. Oncology 1999 17:958-967, Podesta et al., Steroids 1996 61:622-626, Regelson et al., Ann. N.Y. Acad. Sci. 1994 719:564, Schmassmann et al., Gastroenterology 1993 104:1171-1181, Schmassmann et al., Hepatology 1990 11:989-996, Schreiber et al., Lancet 353:459-461, Schreiber, Neth. J. Med. 1998 53:S24-31, Schwartz et al., Cancer Res. 1988 48:4817, Shahidi et al., Biochem. Biophys. Res. Commun. 1999 254:559-565, Steer et al., Ann. Rheum. Dis. 1998 57:732-737, Suzuki et al., Steroids 1998 63:672-677, Suzuki et al., Steroids 1996 61:296-301, Swaan et al., Bioconjugate Chem. 1997 8:520-525, Tang et al, Anticancer Drug Res. 1998 13:815-824, Thomas et al., J. Steroid Biochem. 1986 25:103-108, Utsumi et al., Cancer Res. 1999 59:377-381, Vanden Heuvel, J. Nutr. 1999 129(2S Suppl.):575S-580S, Wang et al., Endocrinology 1998 139:3903-3912, Wong et al., J. Biol. Chem. 1999 274:5443-5453, Xie et al., Endocrinology 1999 140:219-227, Yen et al., Lipids 1977 12:409-413, Zackheim et al., Arch. Dermatology 1998 134:949-954, Zhang et al., Biochim. Biophys. Acta 1991 1096:179-186, Zhu et al., Carcinogenesis 1988 19:2101-2106.\nSome proteins such as interleukin-6 (“IL-6”), erythropoietin (“EPO”) and thrombopoietin (“TPO”) have been examined for their capacity to enhance various aspects hematopoiesis, e.g., Hematology—Basic Principles and Practice, 3rd edition, R. Hoffman, E. J. Benz Jr. et al., editors, Churchill Livingstone, N.Y., 2000 (see, e.g., Chapter 14 at pages 154-202), O. J. Borge et al., Blood 1996 88:2859-2870, M. Cremer et al., Ann. Hematol. 1999 78:401-407, Y. Sasaki et al., Blood 1999 94:1952-1960, U.S. Pat. No. 5,879,673. Recombinant IL-6 was shown in model systems to affect platelet counts in peripheral circulation, e.g., Stahl et al., Blood 1991 78:1467-1475, although significant toxicities are associated with its administration to humans, e.g., Andus et al., FEBS Lett. 1987 221:18, J. Gauldie et al., P.N.A.S. U.S.A. 1987 84:7251-7255, T. Geiger et al., Eur. J. Immunol. 1988 18:717-721. The IL-6 molecule has been described in detail, e.g., publication no. WO 88/00206. Administration of proteins is typically expensive, given factors such as the complexity of producing pharmaceutical grade material.\nThere is a current need for cost-effective pharmaceutical agents and treatment methods for treating various immune dysregulation conditions. The invention provides compounds that can be used in such treatments to treat or ameliorate one or more aspects of immune dysregulation conditions. The compounds can provide surprisingly unexpected beneficial effects, e.g., in preventing or reducing neutropenia in subjects such as mammals or primates. Such agents can be used to treat autoimmune or inflammation conditions, immune suppression conditions, infections, blood cell deficiencies and other described conditions. The agents and methods are useful to ameliorate these conditions or one or more symptoms associated with any of these conditions. The use of these agents can be combined with one or more conventional treatments for these disorders."}
{"text": "The disclosed subject matter is directed to minimizing the exposure to electromagnetic waves and non-ionizing radiation due to use of mobile appliances. This disclosure uses mobile device, mobile appliance, cell phone, cellular phone, wireless phone, wireless device, mobile phone and mobile radio, as well as phone, telephone, cell and radio interchangeably and the use of such terms is not exclusive and is intended to encompass all communication devices and apparatuses that transmit radio waves during communications and are held or located in close proximity to the user.\nMicrowaves are absorbed by living tissues at 24 times the rate of their absorption by pure water. While the Specific Absorption rate, one standard by which cellular microwave absorption by the body is commonly measure, is typically based upon a penetration through an inert emulation of a human head, competent experimental work indicates that the level of absorption in living tissue is many times greater than the level of microwave absorption through an inert liquid, such as water.\nIn a 1983 article entitled An Optical Method for Investigating the Microwave Characteristics of DNA and other Biomolecules in Solution, by Mays L. Swicord and Christopher C. Davis, published in the journal Bioelectronmagnetics, experimental results were reported which determined that “A significant increase in the absorption of DNA solutions as compared with pure water has been observed that is consistent with microwave absorption by the longitudinal mode of the double helix”. Thus, in their conclusion Swicord and Davis state that: However, the results presented in this work are in good agreement with the Prohofsky model of acoustic mode absorption by varying lengths of DNA. Prohofsky and Van Zandt predicted that 450 to 2000 base pair segments of synthetic DNA should absorb 103 to 104 times as strongly in the microwave region as an equivalent mass of water with a decrease in peak absorption due to water damping. The 1.7% dilution of DNA investigated by PFLOH spectroscopy in this work indicated a 40% increase in absorption above pure water at 8 GHz and at 10 to 12% increase at 12 GHz. The measured DC conductance of this DNA sample was quite low yet its absorption coefficient was still 25% higher at 8 GHz than a saline solution of 20 times greater DC conductance. We conclude therefore that the observed absorption of the DNA solution does not come from ionic behavior.\nThe observed absorption is suggestive of direct microwave absorption by the longitudinal acoustic mode of the double helix discussed by Prohofsky and co-workers. Based on the concentration of DNA solution which gave 40% more absorption than pure water at 8 GHz, the microwave absorption of DNA at this frequency is 24 times greater than an equivalent mass of water.”\nThis finding that there are many orders of magnitude of increase in microwave absorption when DNA is present, may likely explain the considerable body of experimental evidence strongly indicating a relationship between long term mobile device use, and diseases involving DNA strand abnormality, such as brain cancer. As, due to the Inverse Square Law, the strength of a radio broadcast or any point source or radial field will be inverse to the square of the distance from the source, exposure levels will diminish rapidly with distance, and a device which occasions increased distance will thereby reduce signal strength absorbed by the human head.\nScientific studies [i.e., H. Lai et al, from the University of Washington, 1984, 1988, and as presented in 1998, Vienna, Austria, and 2009 in Stavanger, Norway; O. Johansson, Associate Professor, Dept. of Neuroscience of the Karolinska Institute, Stockholm, and Professor, Royal Institute of Technology, Stockholm, as presented in 2009 at Stavanger, Norway; Carl F. Blackman a founder of the Bioelectromagnetics Society, as presented in 2009, at Stavanger, Norway; Martin Blank, Ph.D., Associate Professor, Columbia University, as presented in 2009 Stavanger, Norway, Franz Adlkofer, M.D., doctorate from the Max Planck Institute for Biochemistry as presented at Stavanger Norway, 2009, also the following presenters at the International EMF Conference 2009 at Stavanger, Norway: Lukas h. Margaritis, Ph.D., Greece; Elihu D. Richter, MD, M.P.H., Israel; Alvaro Agusto A. de Salles, Ph.D., Brazil; Fredrik Soderqvist, Ph.D., Sweden, Yuri G. Grigoriev M.D. Sci., Russia; Anton V. Merkulov Ph.D., Russia] have shown potential health risks, in some instances showing DNA breakage, as associated with human exposure to non-ionizing radio wave sources, including but not limited to those emitted from mobile telephone devices and handsets, including but not limited to cellular telephones, satellite telephones, cordless telephones, and also portable computers, including those equipped with wi fi connectivity capacity. Continuing research supports that such health risks as may be associated with proximity to non-ionizing radiation will be reduced if level of exposure is reduced.\nAs used here, the term “non-ionizing radiation,” is in reference to types and frequencies of radiation which do not have the capacity to force electron shifting at an atomic level, including with molecular effect. As used here, the term “ionizing radiation” is in reference to types and frequencies of radiation which do have the capacity to force electron shifting at an atomic level, including with molecular effect.\nCurrent electromagnetic field generating communication devices allow the transmission of electromagnetic field though the human head, and are typically held in direct contact with the head of the user. While the SAR standards currently viewed as “safe,” by the FCC rely upon an assumption that the device involved is held at an inch from the head of the user, the practical experience of the typical end user is to hold the device in direct contact with the ear. The inch presumption, however, may be different for several manufactures, but generally lies within ½ to 1¼ inches. Microwave devices, including variants of mobile and portable telephone devices broadcast electromagnetic waves, including but not limited to the radio waves which serve as conduit for the transmission of information, by modulation of amperage or frequency or varying combinations thereof. As discussed above there is competent scientific data indicating that information carrying microwave broadcast near the human body will result in DNA damage and can result in the formation of cancers. In addition, other data indicates that cellular device users with smaller skulls may receive a microwave dosage in inverse relationship to relative skull size, and that, further, the shape of the interior of the skull, including as to concavity, may have a focusing effect, potentially resulting in “hot spots” of electromagnetic signal absorption.\nSome of these deleterious effects have begun to be recognized and have led to some prior art solutions. These prior art solutions however are in contrast to the solutions described in the current application in that the prior art solutions are directed to controlling the emitted signal power level of the transmitter to reduce exposure levels. These prior art solutions include:\nU.S. Pat. No. 7,499,722 to McDowell et al entitled “Proximity Regulation system for Use with a Portable Cell Phone and a Method of Operation Thereof” discloses a portable cell phone including a sensor that determines the proximity of the sensor to a designated part of a user's body (including the head), and controls the transmit power level in response. Thus the transmit power level of the phone is controlled to various “allowable proximity transmit power levels” depending on the vicinity of the portable cell phone to different parts of the user's body.\nU.S. Pat. No. 7,146,139 to Nevermann entitled “System and Method for Reducing SAR levels'” similarly discloses a mobile communications device having a sensor assembly for detecting proximity to a human body and a controller to reduce the average power level of the communication device to less than a predetermined acceptable level. Nevermann discloses that in some embodiments, the power level may be reduced only to the point necessary to maintain adequate signal strength, while in other embodiments the power level may be reduced as required regardless of signal strength in which case a communications link may be dropped.\nU.S. Pat. No. 6,934,515 to Wallach entitled “Electromagnetic Radiation Alerting Device for Use With a Cellular Telephone” discloses a phone that is operable in different modes selected by the user (such as “near”, “normal” and “far” that are calibrated for different distances of the user's head from the antenna), a RF sampling unit for measuring the EMF generated by the cell phone's antenna during its active transmission times and an alerting unit (audio, buzzer, visual) that alerts the user when a predetermined radiation level is exceeded for the selected mode of operation.\nU.S. Pat. No. 6,456,856 entitled “Mobile Radio Equipment forming Antenna Pattern to Protect User From Radiation”. This reference discloses a proximity detection device for measuring humidity or temperature and to detect the presence of human tissue in the vicinity of the device and to adjust transmit power in a directional antenna as a function of the data produced by the detector.\nThe prior art however fails to recognize that providing information to the user in the form of a warning or other alert may be beneficial in maintaining a minimum distance between the user and the mobile device. Without manifestly causing a change in the habits or manner of use of the mobile devices by the user, the prior art approach fails to adequately address the problem.\nIn order to address the detrimental effects described above and obviate the deficiencies of the prior art, the present subject matter discloses controlling the distance between the user, specifically the head, and the mobile device to a minimum distance to therefore reduce exposure. The advantageous reduction in exposure is realized due to the operation of the Inverse Square Rule on electromagnetic signal from cellular and other radio devices, in that there is a steep fall off of signal strength resultant from increased distance between the head of the user and the source of the signal, as a class, devices which result in increased operational distance between the user and the broadcast device will result in a diminution of signal absorption which, while situational variable, can be dramatic in practical effect. As described herein, the disclosed embodiment control the distance between the emitting devices and the user by influencing the user to maintain a minimum distance via alerts, interruptions and warnings.\nThese and many other advantages of the present subject matter will be readily apparent to one skilled in the art to which the invention pertains from a perusal of the claims, the appended drawings, and the following detailed description of preferred embodiments."}
{"text": "Oil spill may have a devastating impact on the surrounding environment. Spilt oil penetrates into the structure of the plumage of birds and the fur of mammals, reducing its insulating ability, and making them more vulnerable to temperature fluctuations and much less buoyant in the water. Clean up and recovery from an oil spill is difficult and may take weeks, months or even years. Therefore, there is a need for improved methods and systems to clean-up oil spills."}
{"text": "Power dissipation associated with networking equipment is increasing as enhanced functionality is designed into various networking products. With enhanced functionality comes an attendant increase in heat load and heat density.\nSome new high-capacity networking products include a bank of optical transceivers known as XFPs (an XFP is a 10 Gigabit/s version of a transceiver device typically known as a small form factor pluggable transceiver or “SFP”) placed directly in line, each dissipating 3.5 W. In normal conditions a commercial-grade SFP/XFP, which can operate up to a maximum of 70° C., is preferred. In more extreme conditions industrial-grade SFPs/XFPs are available, which can operate up to a maximum temperature of 85° C. However, industrial-grade SFPs/XFPs are substantially more expensive than commercial ones.\nPast solutions to cooling a bank of SFPs or XFPs have included an individual heat sink attached to each SFP or XFP. The problem with this design is that SFPs/XFPs have been observed to operate at temperatures exceeding 85° C.\nIn addition, heat sinks commonly have fins which may be of the pin-fin design (round vertical fins) to allow for adequate cooling when the air flowing past the heat sink is expected to be in any orientation. The pin-fin design is suboptimal in heat transfer performance relative to parallel-fin heat sinks, which is a detriment of existing solutions. However, in some applications, the use of pin-fin heat sinks is necessary as some products may be operated in different orientations, such as horizontal and vertical modes.\nA typical bank of SFPs/XFPs is located on a circuit pack and has limited space and height available in the slot occupied by the circuit pack. A typical circuit pack slot is 25 mm in height. The printed circuit board, backside wiring and components and tolerances can account for 7 to 11 mm of this height. As a typical XFP is 10 mm in height, this leaves 4 to 8 mm for the heat sink, which typically includes fins and the heat sink base. As a consequence, the height of the fins on the heat sink is 7-8 mm at most. This limits the amount of heat that the heat sink can pass to the air. For example, typical pin-fin heat sink assemblies that are 4.2 mm and 6.5 mm tall have manufacture-quoted minimum thermal resistances of 18 C/W and 11.5 C/W and are effective at dissipating heat loads of 1.5 W and 2.5 W, respectively.\nSimilar thermal challenges are presented by other circuit pack components. For example, networking products include arrays of three or more FBGAs (fine pitch ball grid arrays) that each dissipate 20 W.\nIn summary, the current solution is deficient as individual heat sinks do not have the ability to adequately dissipate heat to the air under the current space and airflow constraints found in typical circuit packs. This is particularly evident for SFPs and XFPs that are located at the downstream (relative to an incoming airflow) side of a board, where they are required to transfer heat to air that has been heated substantially by the SFPs/XFPs upstream of these components.\nIt is clear from the above that there is a substantial need for a more efficient means of transporting heat away from networking components (e.g. SFPs, XFPs and other heat generating components) and dissipating it to the ambient air."}
{"text": "The present invention relates to a musician's chair, and more particularly to such a musician's chair that can be conveniently and quickly adjusted to the desired elevation and firmly secured against vibration.\nA prior art musician's seat, as shown in FIG. 1, is generally comprised of a stand having a top barrel, a locating ring spaced below the top barrel, and a sleeve connected between the top barrel and the locating ring, a seat with a threaded seat stem threaded into a connector at the top of the sleeve of the stand, and a seat stem lock controlled to lock the seat stem at the desired elevation. The seat stem lock, as shown in FIG. 1, comprises a shell threaded onto the threaded seat stem of the seat and covered on the connector at the top of the sleeve, and a tightening up screw threaded into a radial screw hole and stopped against the periphery of the connector to fix the shell in place. When adjusting the elevation of the seat, the tightening up screw is loosened, and then the seat is rotated upwards or downwards to the desired elevation. Because the threaded seat stem of the seat cannot be moved axially linearly in the connector and the shell, it takes much time to adjust the elevation of the seat."}
{"text": "Automated aerial vehicles, sometimes referred to as drones or unmanned aerial vehicles (UAVs), have become commonly used by hobbyists, some commercial entities, and various government entities. Many of these aerial vehicles are used for image capture, for example, by hobbyists; however, many other uses exist.\nAircraft are often designed for a specific type of flight. Fixed wing aircraft, which primarily rely on a wing for upward lift, typically include a propulsion system that is in fixed orientation and provides thrust in the single direction that, during horizontal flight, is approximately parallel with the horizon. Aircraft that rely on rotors for propulsion, such as helicopters, quadcopters, and other rotorcraft, primarily rely on the rotors for upward lift and typically include a propulsion system that is in fixed orientation and provides thrust in the single direction that, during flight, is approximately perpendicular with the horizon. Some hybrid aircraft rely on both rotors and wings to provide upward lift, depending on a mode of flight. A classic example is the Boeing V-22 Osprey, which uses tilt-rotors that rotate relative to the fuselage to enable transition from a rotor-lift-based mode of flight to a wing-lift-based mode of flight.\nUnmanned aerial vehicles offer unique advantages and considerations as compared to their counterpart manned aerial vehicles (e.g., typical helicopters and fixed wing aircraft). For example, unmanned aerial vehicles may be smaller in overall size and lightweight as compared to their counterpart manned aerial vehicles. The size and thrust ratios of propulsion systems for unmanned aerial vehicles allow these aerial vehicles to perform maneuvers that may not be possible or may not be practical for their counterpart manned aerial vehicles."}
{"text": "In the majority of television viewing environments today, when a user tunes to a specific channel, the graphical interface for that channel is defined primarily by the television manufacturer. For example, when a broadcaster displays a station identifying logo with their video transmission, they have no control over the positioning of their logo. Many times this logo is covered by a user selected menu option when overlaid on the television screen with the transmitted video.\nIn recent environments, graphical interfaces have provided a greater level of functionality with the advent of direct satellite system set top boxes (STB), such as DSS and Echostar. In addition to television displays, the user has channel banners, menus, electronic program guides and message alerts. In some circumstances, such as with inclement weather announcements and breaking news events, the user may customize the interface. However, the amount of functionality presented by these graphical interfaces is minimal since there is little interaction by a user, other than changing the channel or adjusting the volume.\nIn particular, these functions provided by an STB manufacturer occupy a portion of the display screen when selected by the user, such that the television broadcaster's logo or call letters are very likely to be covered up if displayed on the display screen at the same time a user is operating the electronic program guide, for example. The broadcasters have no control over the positioning of their station identifying information or other types of data they may wish to have in full view when displayed concurrently with user selected display options. With the DSS and Echostar systems, the additional user interfaces have been designed and implemented by STB manufacturers with little regard for the needs or concerns of the television broadcaster.\nTherefore, there is a direct conflict between the manufacturer's need to define the various functions provided by a direct satellite STB and the broadcaster's need to broadcast unobstructed graphical symbols and other information to viewers.\nFurthermore, this problem is not just limited to television broadcasts and receptions within the United States, but worldwide. This tension between the manufacturers and the broadcasters will increase as the boundaries of televisions and computing devices blur with the evolution of televisions adapting to the digital age and as data services become more readily available.\nThis convergence between televisions and computers can be illustrated by the efforts of Microsoft and Intel, where efforts have been made in developing strategies and approaches for integrating televisions into the personal computer (PC). For instance, Broadcast PC has been developed by Microsoft and Intel has developed Intercast, but they do not allow for the broadcaster's total control of the television environment. In both developments, a broadcaster is limited to controlling the look and feel of their services, but has no control of their primary product, which is the broadcast event.\nWindows CE, OS-9 and OpenTV are all operating systems that are attempting to deliver more sophisticated user interfaces to the manufacturer, but they too are disregarding the needs of the broadcaster. These operating systems still cover up the broadcast content with banners and menus that are genericized to the lowest common denominator of services and functions offered.\nTherefore, there is a need to permit broadcasters to have greater control in the television environment they transmit to their viewers. Specifically, there is a need for the broadcasters to broadcast unobstructed graphical symbols and other information to viewers."}
{"text": "As the supply of low sulfur, low nitrogen crudes decrease, refineries are processing crudes with greater sulfur and nitrogen contents at the same time that environmental regulations are mandating lower levels of these heteroatoms in products. Consequently, a need exists for increasingly efficient desulfurization and denitrogenation catalysts.\nIn one approach, a family of compounds, related to hydrotalcites, e.g., ammonium nickel molybdates, has been prepared. Whereas X-ray diffraction analysis has shown that hydrotalcites are composed of layered phases with positively charged sheets and exchangeable anions located in the galleries between the sheets, the related ammonium nickel molybdate phase has molybdate anions in interlayer galleries bonded to nickel oxyhydroxide sheets. See, for example, Levin, D., Soled, S. L., and Ying, J. Y., Crystal Structure of an Ammonium Nickel Molybdate prepared by Chemical Precipitation, Inorganic Chemistry, Vol. 35, No. 14, p. 4191-4197 (1996). The preparation of such materials also has been reported by Teichner and Astier, Appl. Catal. 72, 321-29 (1991); Ann. Chim. Fr. 12, 337-43 (1987), and C. R. Acad. Sci. 304 (II), #11, 563-6 (1987) and Mazzocchia, Solid State Ionics, 63-65 (1993) 731-35.\nNow, when molybdenum is partially substituted for by tungsten, an amorphous phase is produced which upon decomposition and, preferably, sulfidation, provides enhanced hydrodenitrogenation (HDN) catalyst activity relative to the unsubstituted (Ni--Mo) phase."}
{"text": "The term amplifier is used below to describe a device which can be used to increase the intensity of an optical beam.\nFIG. 1 shows a segment of a section of optical waveguide. An optical waveguide L is connected via an optical connector S to an optical amplifier V. Optimum noise behavior of an optical amplifier, in particular an erbium doped fiber amplifier, results if the amplifier does not have isolators on the input side. Amplified spontaneous emission ASE caused by the optical amplifier spreads counter to the direction of spread of the optical signal on the optical waveguide and is, for example, reflected via a connector S and fed back to the input of the amplifier.\nIn order to achieve safety shutdown of the optical amplifier in the event of interruption, components of the optical signal directed to the input of the optical amplifier V can be fed to a monitoring unit via an optical coupler in appropriate transfer systems. This monitoring unit is, for example, made up of a filter, particularly a band-pass filter, and a signal monitor downstream from this. If the measured power is below a predefined threshold value, it is assumed that the optical waveguide section is interrupted, which, for example, triggers safety shutdown of the optical amplifier. If the measured power is above the predefined threshold value, safety shutdown does not take place, as it is assumed that optical signals have reached the measuring device via the optical waveguide section and, therefore, there can be no interruption.\nIn the case of interruption, however, a power reading above the predefined threshold value may be measured under certain circumstances due to reflection of the amplified spontaneous emission ASE at a connector S at the input of the amplifier and this is interpreted as the input signal, so the mandatory safety shutdown does not take place.\nTo ensure shutdown in the event of an interruption, optical amplifiers are therefore fitted with isolators on the input side.\nThis has the advantage that amplified spontaneous emission ASE spreading counter to the direction of spread of the optical data signal is avoided.\nThe use of an isolator, however, such as a passive element, has the disadvantage that the noise level increases and, as a result, the intervals between the regenerators along the optical transfer section are reduced.\nAn object of the present invention is to specify a circuit and a method for safety shutdown of an optical amplifier."}
{"text": "Low-emissivity coatings are well known in the art. Typically, they include one or two layers of infrared-reflection film and two or more layers of transparent dielectric film. The infrared-reflection film, which generally is a conductive metal like silver, gold, or copper, reduces the transmission of heat through the coating. The dielectric films are used to anti-reflect the infrared-reflection films in selected spectral regions (visible region of sun light) and to control other properties and characteristics of the coating, such as color and durability. Commonly used dielectric materials include oxides of zinc, tin, indium, bismuth, and titanium, among others.\nMost commercially available low-emissivity coatings have one or two silver layers each sandwiched between two coats of transparent dielectric film. Increasing the number of silver films in a low-emissivity coating can increase its infrared reflection. In commercial coating processes the Ag films are supported by growth support layers (below) and protective layers above. Especially the protective layers above can add some additional (to the silver films) absorption to the layer-stack. For multiple silver films in a layer stack this can add up to total absorption levels not acceptable for some applications. This will also reduce the visible transmission of the coating, and/or negatively impact the color of the coating, and/or decrease the durability of the coating. In some processes, growth support layers are provided beneath silver films and protective layers are provided above silver films. An increase in the number of growth support layers and protective layers in a low-emissivity coating can increase the overall absorption of the coating. This can be undesirable in some cases. Perhaps for these reasons, low-emissivity coatings with three silver layers have not found much place in the market.\nIt would be desirable to provide a low-emissivity coating that includes three infrared-reflection film regions and has desirable coating properties and characteristics. It would also be desirable to provide deposition methods and equipment that can produce high quality coatings of this nature."}
{"text": "Demand for 1,3-butadiene, as a petrochemical intermediate, and value thereof are gradually increasing.\nAs methods of preparing 1,3-butadiene, there are naphtha cracking, direct dehydrogenation of butene, oxidative dehydrogenation of butene, etc.\nThereamong, oxidative dehydrogenation of butene, as a reaction generating 1,3-butadiene and water by reaction of butene with oxygen, generates stable water, thus being very thermodynamically advantageous.\nIn addition, since the oxidative dehydrogenation of butene is an exothermic reaction unlike the direct dehydrogenation of butene, it enables production of 1,3-butadiene in a high yield even at low reaction temperature, compared to direct dehydrogenation. In addition, in the case of the oxidative dehydrogenation, additional heat supply is not necessary, thereby being very suitable for a commercialization process.\nHowever, after the oxidative dehydrogenation of butene, the intensity of a catalyst is decreased, thereby elevating differential pressure of a reactor."}
{"text": "In an earlier form, an articulated, predominantly concrete mat is disclosed in U.S. Pat. No. 5,108,222, the disclosure of which is incorporated by reference herein. In an improved form, an articulated, predominantly concrete mat is disclosed in U.S. Pat. No. 6,612,776 B1, the disclosure of which is incorporated by reference herein. As disclosed in these patents, such a mat has numerous uses in retarding earth erosion due to wind, water, or both and in lining a driveway, a parking area, or an emergency roadway.\nU.S. Pat. No. 6,612,776 B1 discloses that the articulated, predominantly concrete mat comprises concrete blocks, which are connected to one another by flexible members of a flexible geogrid. As disclosed therein, the articulated, predominantly concrete mat is made from a concrete slab, in which the flexible geogrid is embedded. As disclosed therein, the concrete slab has relatively thicker portions, which define such blocks, and relatively thinner portions, along which the concrete slab is broken to form such blocks and which has holes to facilitate breaking of the concrete slab. As disclosed therein, the flexible members of the flexible geogrid are flexible straps, as disclosed in U.S. Pat. No. 5,108,222, or are flexible strands or strand bundles.\nTypically, an articulated, predominantly concrete mat, as discussed above, is sized to cover a ground area very many times smaller than a flexible geogrid, as available commercially, is able to cover. Thus, before being employed in an articulated, predominantly concrete mat, as discussed above, the flexible geogrid must be first cut to a smaller size. In a flexible geogrid, as available commercially, if spacings between the flexible members are imprecise in a longitudinal direction, in a lateral direction, or in both directions, it is difficult to employ the flexible geogrid in an articulated, predominantly concrete mat, as discussed above.\nMoreover, in an articulated, predominantly concrete mat, the flexible members of the flexible geogrid reinforce the concrete blocks, except that the concrete blocks along the longitudinal and lateral edges of the articulated predominantly concrete mat may not be sufficiently reinforced by the flexible members of the flexible geogrid to prevent outer corners of the concrete blocks from cracking or breaking.\nFurthermore, when a flexible geogrid, as available commercially, is embedded in a concrete slab, which is cast from a concrete slurry, the flexible geogrid tends to float upwardly in the concrete slurry, before the concrete mat has cured.\nFor any one or more of the foregoing reasons, a flexible geogrid, as available commercially, may prove to present shortcomings disfavoring its use in an articulated, predominantly concrete mat, as discussed above."}
{"text": "With the multitude of operating systems and versions thereof, developers have created software that provides virtual machines (sometimes referred to as virtual PCs). A virtual PC is a software-created environment that may appear, to the operating system and applications that execute therein, to have most or all the elements associated with an actual PC. The software-created environment may be so complete that the operating system and applications that execute on a virtual PC may not be able to determine that they are executing in a virtual environment.\nProviding virtual PCs is useful for many reasons including debugging and supporting applications designed to execute on multiple versions of one or more operating systems, having access on a single machine to applications that are only available for different operating systems, and so forth. To execute multiple operating systems on one actual computer, a user may load a primary operating system on the computer, load software to create one or more virtual PCs, and then load one or more other operating systems that execute in the one or more virtual PCs. This allows, for example, an Apple® (or other) operating system to provide a virtual PC to host a Microsoft® operating system, a UNIX® operating system, a LINUX® operating system, another Apple® operating system, or some other operating system.\nProviding a virtual PC that seamlessly and effectively provides all the functionality of a real PC has eluded those skilled in the art. This is particularly true with respect to printing functions associated with the virtual PC."}
{"text": "Typically, in distributed server environments, many different physical servers contain duplicate information, thereby permitting large numbers of clients to access the same information from different servers. Often, this is accomplished by directing each query to a different server in a round-robin fashion. This type of distributed server environment provides faster access to data, since much of the processing and network traffic is distributed to many different servers.\nGenerally, the distributed servers are updated by pushing information to each of the servers. In other words, if the contents of the servers need to be updated, then the new contents are often uploaded to each of the servers from a central repository. Unfortunately, when a vast number of servers need to be updated concurrently, the bandwidth is occupied during the time of update, thereby slowing network traffic. In view of this deficiency, a need exists in the industry for controlling and maintaining information in distributed server environments."}
{"text": "1. Field of the Invention\nThe disclosure relates to a nonvolatile memory device and a method of manufacturing the same, and more particularly, to a nonvolatile memory device having an increased trap site density and a method of manufacturing the same.\n2. Description of the Related Art\nIn semiconductor memory devices, especially DRAMs (Dynamic Random Access Memory), a unit memory cell includes one transistor and one capacitor. Therefore, to increase the integration density of semiconductor memory devices, the volume of the transistor or the capacitor must be reduced.\nIn early semiconductor memory devices when the integration density was not a large concern, photography and etching processes could be performed with sufficient process margins. The integration density of semiconductor memory devices could be increased by reducing the volume of each element of the memory device.\nHowever, as the demand for highly integrated semiconductor memory devices increases, a different method to reduce volume is needed.\nThe integration density of a semiconductor memory device is closely related to a design rule. Therefore, to increase the integration density, the design rule must be strict. In this case, the process margins of photography and etching processes can be significantly reduced, meaning that these processes must be performed more precisely.\nWhen the process margins of the photography and etching processes are reduced, the yield can also be reduced. Therefore, a method to increase the integration density of semiconductor memory devices without reducing the yield is needed.\nAccording to this requirement, many semiconductor memory devices having different structures than conventional semiconductor memory devices have been introduced by including a data storing medium that can store charge on the upper side of a transistor, with a different data storing function.\nThe SONOS (Semiconductor-Oxide-Nitride-Oxide-Semiconductor) memory device is another newly introduced semiconductor memory device. FIG. 1 is a cross-sectional view illustrating a conventional memory device.\nReferring to FIG. 1, a source region 12 and a drain region 14 to which an n type conductive dopant is implanted on a p type semiconductor substrate 10 (hereinafter, semiconductor substrate) are formed. A channel region 16 is formed between the source region 12 and the drain region 14. Also, a gate stack 18 is formed on the channel region 16 of the semiconductor substrate 10. The gate stack 18 is composed of a tunneling oxide film 18a, a nitride film Si3N4 18b, a blocking oxide film 18c, and a gate electrode 18d. Here, the nitride film 18b has trap sites. Therefore, when a voltage is applied to the gate electrode 18d, electrons pass through the tunneling oxide film 18a and are trapped in the trap site of the nitride film 18b. The blocking oxide film 18c blocks the migration of electrons to the gate electrode 18d while the electrons are trapped.\nIn this conventional semiconductor memory device, binary scale information can be stored and read using the characteristic that the threshold voltage varies depending on whether electrons are trapped in the trap site of the nitride film 18b. \nWhen the density of a trap site increases, more electrons can be trapped, and the variation of the threshold voltage can be increased. That is, the density of the trap site can significantly affect the characteristics of the memory device. Conventionally, to increase the density of a trap site, techniques of scattering or depositing nano-scale particles on the surface of a thin film have been developed. However, these methods can only provide limited increases in the density of the trap site per unit area. These methods have various technical problems, especially with respect to uniformity in the flash memory. Also, in the case of conventional memory devices, the only way to increase the integration density is to reduce the volume. However, since this makes the design rule more strict, there is a limit to increasing the integration density of memory devices by reducing their volume."}
{"text": "Dielectric resonator antennas are resonant antenna devices that radiate or receive radio waves at a chosen frequency of transmission and reception, as used in for example in mobile telecommunications. In general, a DRA consists of a volume of a dielectric material (the dielectric resonator or pellet) disposed on or close to a grounded substrate, with energy being transferred to and from the dielectric material by way of monopole probes inserted into the dielectric material or by way of monopole aperture feeds provided in the grounded substrate (an aperture feed is a discontinuity, generally rectangular in shape, although oval, oblong, trapezoidal or butterfly/bow tie shapes and combinations of these shapes may also be appropriate, provided in the grounded substrate where this is covered by the dielectric material. The aperture feed may be excited by a strip feed in the form of a microstrip transmission line, coplanar waveguide, slotline or the like which is located on a side of the grounded substrate remote from the dielectric material). Direct connection to and excitation by a microstrip transmission line is also possible. Alternatively, dipole probes may be inserted into the dielectric material, in which case a grounded substrate is not required. By providing multiple feeds and exciting these sequentially or in various combinations, a continuously or incrementally steerable beam or beams may be formed, as discussed for example in the present applicant's co-pending U.S. patent application Ser. No. 09/431,548 and the publication by KINGSLEY, S. P. and O'KEEFE, S. G., “Beam steering and monopulse processing of probe-fed dielectric resonator antennas”, IEEE Proceedings—Radar Sonar and Navigation, 146, 3, 121–125, 1999, the full contents of which are hereby incorporated into the present application by reference.\nThe resonant characteristics of a DRA depend, inter alia, upon the shape and size of the volume of dielectric material and also on the shape, size and position of the feeds thereto. It is to be appreciated that in a DRA, it is the dielectric material that resonates when excited by the feed. This is to be contrasted with a dielectrically loaded antenna, in which a traditional conductive radiating element is encased in a dielectric material that modifies the resonance characteristics of the radiating element.\nDRAs may take various forms, a common form having a cylindrical shape dielectric pellet which may be fed by a metallic probe within the cylinder. Such a cylindrical resonating medium can be made from several candidate materials including ceramic dielectrics.\nSince the first systematic study of dielectric resonator antennas (DRAs) in 1983 [LONG, S. A., McALLISTER, M. W., and SHEN, L. C.: “The Resonant Cylindrical Dielectric Cavity Antenna”, IEEE Transactions on Antennas and Propagation, AP-31, 1983, pp 406–412], interest has grown in their radiation patterns because of their high radiation efficiency, good match to most commonly used transmission lines and small physical size [MONGIA, R. K. and BHARTIA, P.: “Dielectric Resonator Antennas—A Review and General Design Relations for Resonant Frequency and Bandwidth”, International Journal of Microwave and Millimetre-Wave Computer-Aided Engineering, 1994, 4,(3), pp 230–247]. A summary of some more recent developments can be found in PETOSA, A., ITTIPIBOON, A., ANTAR, Y. M. M., ROSCOE, D., and CUHACI, M.: “Recent advances in Dielectric-Resonator Antenna Technology”, IEEE Antennas and Propagation Magazine, 1998, 40, (3), pp 35–48.\nA variety of basic shapes have been found to act as good DRA resonator structures when mounted on or close to a ground plane (grounded substrate) and excited by an appropriate method. Perhaps the best known of these geometries are:\nRectangle [McALLISTER, M. W., LONG, S. A. and CONWAY G. L.: “Rectangular Dielectric Resonator Antenna”, Electronics Letters, 1983, 19, (6), pp 218–219].\nTriangle [ITTIPIBOON, A., MONGIA, R. K., ANTAR, Y. M. M., BHARTIA, P. and CUHACI, M.: “Aperture Fed Rectangular and Triangular Dielectric Resonators for use as Magnetic Dipole Antennas”, Electronics Letters, 1993, 29, (23), pp 2001–2002].\nHemisphere [LEUNG, K. W.: “Simple results for conformal-strip excited hemispherical dielectric resonator antenna”, Electronics Letters, 2000, 36, (11)].\nCylinder [LONG, S. A., McALLISTER, M. W., and SHEN, L. C.: “The Resonant Cylindrical Dielectric Cavity Antenna”, IEEE Transactions on Antennas and Propagation, AP-31, 1983, pp 406–412].\nHalf-split cylinder (half a cylinder mounted vertically on a ground plane) [MONGIA, R. K., ITTIPIBOON, A., ANTAR, Y. M. M., BHARTIA, P. and CUHACI, M: “A Half-Split Cylindrical Dielectric Resonator Antenna Using Slot-Coupling”, IEEE Microwave and guided Wave Letters, 1993, Vol. 3, No. 2, pp 38–39].\nSome of these antenna designs have also been divided into sectors. For example, a cylindrical DRA can be halved [TAM, M. T. K. and MURCH, R. D.: “Half volume dielectric resonator antenna designs”, Electronics Letters, 1997, 33, (23), pp 1914–1916]. However, dividing an antenna in half, or sectorising it further, does not change the basic geometry from cylindrical, rectangular, etc.\nHigh dielectric antennas (HDAs) are similar to DRAs, but instead of having a full ground plane located under the dielectric pellet, HDAs have a smaller ground plane or no ground plane at all. Removal of the ground plane underneath gives a less well-defined resonance and consequently a very much broader bandwidth. HDAs generally radiate as much power in a backward direction as they do in a forward direction.\nIn both DRAs and HDAs, the primary radiator is the dielectric pellet. In DLAs, the primary radiator is a conductive component (e.g. a metal wire or printed strip or the like), and a dielectric component then just modifies the medium in which the DLA operates and generally allows the antenna as a whole to be made smaller or more compact."}
{"text": "1. Field of the Invention\nThe present invention is directed to a heat sink fastening device for attaching a heat sink to an integrated circuit and a socket receiving the integrated circuit.\n2. Description of the Prior Art\nThe overheating of integrated circuits and central processing units (CPU), in particular, has been the critical issue that needs to be seriously addressed for a computer system to function normally without system failure. This is especially true in the present day when the clock speed of a typical CPU is so fast and the requisite computations are so great that an increased amount of heat is generated from it. In order to carry the waste heat away from the CPU and therefore lower its operating temperature, a heat sink device is conventionally utilized to effect the dissipation of heat.\nFIG. 1 illustrates how a heat sink is clamped onto a CPU according to the prior art. A CPU 103 is to be inserted into a socket 104 specially designed for this purpose on a motherboard, for example. A heat sink 102 with arrays of fin structure is attached to the upper surface of the CPU 103. The contact between the heat sink 102 and the CPU 103 is supposed to be tight and the contact surface as large as possible, so that the heat transfer between them can be optimized. For such purpose, a clip 100 is used to tightly and securely fasten the heat sink 102 upon the CPU 103. The clip 100 has an elongated elastic main body that enters into the spacing 105 of the heat sink 102 and contacts the base surface of the heat sink 102 at a pressing portion 107 on the main body. At the opposite ends of the main body are two latching portions 101 each having an opening for engaging with two protrusions 106 disposed on opposite sides of the socket 104, and thereby tightly clamping the heat sink 102, the CPU 103, and the socket 104 altogether.\nUsually the clip 100 includes a lever arm 108 to facilitate mounting and removing of the clip 100. The clamping force is derived from the M shape of the clip 100, as is illustrated in FIG. 2, which deforms elastically when the latching portions 101 are locked onto the protrusions 106 and the pressing portion 107 is contacting and pressing the heat sink 102 against the CPU 103. Although the clip 100 is capable of exerting a downward force upon the heat sink 102 when mounted thereon due to elastic deformation, it simultaneously exerts an outward force upon the protrusions 106. This is because that, when mounting the clip 100, the latching portions 101 have to be pushed inwardly towards each other in order to lock onto the protrusions 106, so that the resilient force of the latching portions 101 in the locked position is directed outward. This is disadvantageous, for the outward resilient force tends to disengage the latching portions 101 out of the protrusions 106, especially during sudden vibration and percussions of the whole assembly."}
{"text": "A montage may comprise an aggregation of separate elements to form a single composite element. For instance, a montage may comprise a composite picture made up of several separate pictures, or a video sequence comprising a rapid sequence of disparate images. Electronic montage systems have been designed to create digital montages using digital content, such as a composite web page comprising different constituent web pages served from different web applications. Sometimes the constituent web pages are organized according to a central theme, such as a web pages related to a given search term used by a search engine, or web pages consistently visited by a user as stored in a browser history. Often, the constituent web pages are lower-fidelity representations of the actual web page due to space limitations of the composite web page. As such, a user may select a constituent web page to retrieve a higher-fidelity version of the selected web page for more in-depth viewing. As an amount of digital information increases, however, it becomes increasingly difficult to build a digital montage in a way that provides meaningful information to a user. It is with respect to these and other considerations that the present improvements have been needed."}
{"text": "1. Field of the Invention\nThe invention relates to magnetic disk drives and, in particular, to a method and apparatus for automatic loading/unloading of a magnetic head in such a magnetic disk drive.\n2. Description of the Prior Art\nIn computer systems, information is frequently stored on a magnetic film on the surface of a magnetic disk (hard disk). Magnetic hard disk storage devices, that is to say assemblies of several such magnetic disks, either have permanently fixed disks or removable cartridges which can be removed from a disk storage device and inserted into the next device of the same type. These disks are not perfectly flat; but due to the manufacturing process, the disks have fine curvatures. Information is read and written by means of magnetic heads (sliders) which are connected to the housing of the disk storage device by means of a carrier structure. The magnetic heads, together with their suspensions, are mounted on the arms of the carrier structure, the actuator.\nFIG. 1 shows a part of such an assembly according to the state of the art. The magnetic head suspensions themselves consist of the loadbeam, the flexure, and the mounting block, collectively referred to as the suspension. The mounting block serves as a connecting element between the actuator arm (not shown) and the magnetic head suspension. The loadbeam is the backbone of this suspension. It has a relatively high rigidity against the slider. Toward the rear, it has a controlled flexibility defined in production. In this area, bending of the loadbeam in z direction furnishes the cross-sectional profile with a defined rigidity. When the actuator is subsequently fitted in the disk storage device, this controlled rigidity of the bend zone in the fitting position produces a pressure load (gramload) of the slider on the surface of the disk.\nThe flexure is welded onto the loadbeam. The slider is then bonded onto a lug of the flexure. Due to its shape, the flexure has a very high lateral rigidity to prevent the slider from oscillating uncontrolledly in the event of rapid actuator movements. The rigidity around the transverse axis, on the other hand, is very low, to allow the slider to flex in response to the minutest unevenness on the disk. The entire assembly is referred to as the Head Suspension Assembly (HSA).\nDuring read and write operations, the slider floats on a cushion of air at a very low height above the disk, which rotates very fast. The bearing surface onto the disk is termed the Air Bearing Surface (ABS). In this way, direct contact between the slider and the disk is avoided. If, for example, the rotational speed of the disk is too low to create an adequate air cushion between the head and the disk, the two may touch, resulting in a \"head crash\", which may cause irreparable damage to the disk and head and may lead to loss of data.\nEspecially in replaceable cartridge systems in which the magnetic disk is inserted into and withdrawn from the storage device housing inside a cartridge system, a mechanism is required to lift the magnetic heads off of the disk surface and set them to a rest position when the drive is not in operation. Otherwise, there would be very great danger of damage, in transit for example. However, a mechanism of this kind is also necessary in fixed systems, to prevent the head sticking to the disk at rest position.\nIn earlier times, it was usual to start and stop a magnetic disk, or magnetic disk storage device, by means of a Contact Start/Stop System (CSS). In this process, the drive is started and stopped while the magnetic disk and head are in contact with each other. To prevent the head and disk sticking together as already mentioned when the drive is at rest, the surface of the disk had to be appropriately roughened. Due to the increasing miniaturization of hard disks and hard disk storage devices, however, the gap between the actuator and the surface of the disk is becoming ever smaller, so that this solution no longer corresponds to the state of the art. As a result, it is becoming more and more problematic to unload the magnetic heads with mechanisms of sufficient strength. Also, reducing the size of gap means that the influence of production and assembly tolerances is becoming ever greater.\nFor this reason, load/unload mechanisms have been developed which lift the head from the disk using a load/unload element (from now on, for the sake of simplicity, termed an unload element) when the disk is not rotating. When the disk is started again, i.e., the slider is brought back onto the disk surface, the appropriate air cushion must be reestablished fast enough to prevent contact between the disk and the head, dependent on the lifting speed of the head, the position of the ABS relative to the disk, and the defined height of flight of the head. Because users of such drives wish to read and write data very shortly after it starts operating, a mechanism of this kind must be furnished with the highest possible load/unload speed.\nWO87/01853 discloses a load/unload mechanism for magnetic heads which has a comb-like structure as the unload element. The fingers of the element interact with the suspension of the magnetic heads to hold the heads at sufficient distance from the disk in the unloaded position. This device has the disadvantage that, due to its size, it can no longer be used for today's much more miniaturized magnetic disk drives.\nU.S. Pat. No. 5,296,986 describes a device to hold the actuator of a disk drive in a secure position when the drive is at rest. For this, the slider is conveyed with the aid of a driver mechanism onto a ramp outside the magnetic disk, where it rests. The disadvantage of this device is that additional means such as cam followers are required to convey the head onto the ramp. Another disadvantage is that abrasion occurs in the contact zone between the ramp surface and the slider suspension. The abrasions very often take on magnetic properties due to changes in the grid structure. At the point at which such magnetic particles touch the disk surface, the data already written to it would be destroyed.\nU.S. Pat. No. 5,394,281 describes a magnetic load/unload mechanism with a ramp as the unload element, having a piezoelectric element. In this mechanism, the static frictional force at the start of loading is replaced by a dynamic force. A magnetic solution of this kind has the disadvantage that the slider may drop out of control onto the surface of the disk, for example, if the power fails. Furthermore, this arrangement is not suitable for loading and unloading several magnetic heads simultaneously. In this case, coupling of several different transducers is not possible; and the available space for attachment of several such transducers is inadequate with today's low design heights.\nA further disadvantage of the general state of the art is that when the suspension conveyed runs onto the unload element mounted outside the disk, the risk of collision between the magnetic head and disk is very great, for the following reason.\nIn production, the magnetic disk cannot be fully tested, i.e., the magnetic and mechanical properties are not precisely defined on the entire disk. Thus, due to technical factors, the outermost zone of the disk remains undefined. The consequence is the outer edge of the disk is much less flat than the zones further in. When the suspension is then moved onto the unload element with the disk rotating, the relative movement of the disk and the suspension causes the magnetic head to skew. This leads to a loss of flying height of the magnetic head and so to the risk of solid body contact between the head and the disk at the outermost, uneven, edge of the disk. For this reason, the suspension is normally unloaded in the outermost defined data field of the disk.\nAccordingly, it can be seen that there is a need for a method and device for automatic loading and unloading of a magnet head over a disk surface in a magnetic disk drive which allows the magnetic heads of the drive to be loaded at the highest possible speed in a secure manner and without collision with the surface of the disk."}
{"text": "1. Field of the Invention\nThe present invention relates to a compound floating magnetic head that writes and reads data to and from a magnetic recording medium such as a magnetic disk, and more particularly to a compound floating magnetic head including sealing glass that is chemically stable even in a highly humid environment.\n2. Prior Art\nA conventional compound floating magnetic head uses dynamic pressure generated by the travel of a magnetic medium such as a magnetic disk to float the slider of the head with a very small gap between the head and the medium. The ferrite core disposed at the slider section where the gap is the narrowest writes and read data.\nThe compound floating magnetic head comprises a ferrite core 20 composed of an I-shaped bar 21 and a C-shaped bar 22 shown in FIG. 1 (a), and a non-magnetic slider 10 shown in FIG. 1 (b).\nThe I-shaped bar 21 and the C-shaped bar 22 are coupled at two positions (front gap FG and back gap BG) and integrated into the ferrite core 20. Front gap FG is positioned on the side of the recording medium (not shown). The length of the gap is determined to a specified width depending on the recording density of the recording medium.\nTo maintain the gap length at the (specified width), reinforcing bonding glass 23 is formed at the sloped portion on the side of front gap FG.\nFurthermore, on the side of front gap FG, the front end of the ferrite core 20 has a step with a specified width that depends on the track width of the recording medium.\nThe slider 10 is basically a rectangular parallelopiped. To float the slider 10 apart from the recording medium with a specified gap, floating rails 12 and 13 and bleed groove 14 are formed on the surface 11 of the slider 10 that faces the recording medium.\nA core accommodation groove (core slit) 16 with a specified depth is formed in the longitudinal direction on the trailing end surface 15 of the floating rail 12.\nAs shown in FIG. 1 (c), the ferrite core 20 is coupled with a non-magnetic slider 10 in the core slit 16 of the slider 10 by using sealing glass 17. An coupling process is explained below referring to FIG. 1 (d). The ferrite core 20 is temporarily held in the core slit 16 of the slider 10. A bare 18 that is to be melted to form the sealing glass 17 is then placed on the core 20. By heat treatment at a specified temperature, the bar 18 is softened and melted to fill the core slit 16 with the sealing glass 17, thus securing the ferrite core 20 to the slider 10.\nIn the above-mentioned coupling process, it is important to set heat treatment temperature T1 required to soften and melt the bar 18 to form the sealing glass 17 sufficiently lower than softening temperature T2 of the reinforcing bonding glass 23 of the ferrite core 20 (T1&lt;T2). If heat treatment temperature T1 is set higher than softening temperature T2 of the reinforcing bonding glass 23, the specified electromagnetic conversion characteristics required between the magnetic medium and the ferrite core 20 are deteriorated. For example, the specified gap length on the front gap (FG) side becomes larger.\nConventionally, glass including much lead oxide was used as the sealing glass 17 to obtain a lower softening point. The sealing glass 17 that was used generally comprises 78% PbO by weight, 2.6% SiO.sub.2 by weight and 10% B.sub.2 O.sub.3 by weight. By including much PbO and B.sub.2 O.sub.3 as described above, the sealing glass 17 was easily able to have a low softening point sufficiently lower than softening point T2 of the bonding glass 23. However, when a conventional compound floating magnetic head with this kind of sealing glass 17 was left in a very humid environment, a moisture absorption layer was formed on the sealing glass 17. Chemically, the metal ions included in the sealing glass 17 reacted with hydroxyl groups, weathering and corroding the glass surface. These caused problems: discoloration of the exterior, damage to the recording medium due to dropping of substances precipitated by chemical reaction, and attachment of the magnetic powder of the recording medium to cavity portions from which the precipitated substances were separated.\nIn any case of the above-mentioned problems, noise was caused during writing and reading data to and from the recording medium, significantly reducing the reliability of the electromagnetic conversion characteristics."}
{"text": "For many years, computers typically utilized only one hardware page size: 4 KB. Many software applications were built on assumptions about this page size and relied on its usage. Today, this hardware page size is still commonly in use (such as in x86 architectures), but for performance reasons, larger page sizes are created at the virtual memory level (i.e., within the operating system or software level, as distinguished from the 4 KB hardware pages allocated within physical memory). For example, some operating systems and databases utilize virtual memory pages which span multiple megabytes or even gigabytes.\nGenerally, an operating system creates a data structure to track each physical memory page utilized by virtual memory. For example, in Linux, each physical memory page frame has a small corresponding data structure (struct page) also stored in memory that is used to track the status of the page frame. These data structures by themselves consume a fixed percentage of the total physical memory in the machine, approximately 1% of a system's physical memory. For example, on a machine with 1TB of memory, about 10 GB of memory is dedicated these data structures. As the amount of physical memory in the system grows, so does the number and size of data structures needed to track the page frames in physical memory. (For simplicity, these data structures tracking physical memory page frames in the kernel are hereinafter referred to as “kernel memory management data structures”).\nOne ongoing trend in computer systems is an increase in page sizes used within an operating system's virtual memory. The use of larger page sizes within virtual memory provides significant advantages for modern software applications. However, the use of larger page sizes is also likely to waste a larger amount of memory with unnecessary kernel memory management data structures, as the number of base hardware pages assigned to each large page increases. For example, with 16 MB page mappings being used on a system with a 4K base page size, this would mean that there are 4096 hardware pages mapped to every 16 MB virtual memory page. Therefore, 4096 kernel memory management structures would be associated with each virtual memory page. In practice, however, only the information of the first 4K hardware page will be utilized to determine the starting point for the 16 MB page if the 16 MB page is located contiguously in 4K chunks in physical memory. With the use of larger virtual memory page sizes, a significant number of the kernel memory data structures tracking the hardware pages—as many as 4095 in the preceding example—would be redundant and unused.\nOne workaround to the problem of wasted memory allocated to kernel memory management structures might be to track every single hardware page in a tree structure, adding entries for pages one at a time. Although this would prevent wasting memory on redundant structures, such a tree would be slow and complex to access. Likewise, if all memory was managed in one single flat structure, it would be simple and fast to access, but would be inflexible and waste space for areas not currently in use. To enable performance by efficiently adding memory in the largest amounts possible (such as when hotplugging memory), enhanced techniques for managing existing kernel memory management structures are needed. The techniques utilized within the presently disclosed invention attempt to solve the above-described limitations by reducing the overhead of these memory structures and the need to store redundant memory structures in physical memory."}
{"text": "1. Technical Field\nThe present invention generally relates to a semiconductor device and, more particularly, to a semiconductor device having a plug for diffusing hydrogen into a semiconductor substrate which is covered by one or more layers which substantially form a barrier to hydrogen diffusion.\n2. Description of Related Art\nJunction leakage is a well-known problem for semiconductor devices such as transistors and can adversely affect device operation and performance. One technique for reducing junction leakage diffuses hydrogen into the semiconductor substrate on which the semiconductor devices are fabricated. This diffusion may be accomplished by an annealing process (typically at a temperature in the range from about 400.degree. to about 600.degree. C.) in an atmosphere containing hydrogen. The source of hydrogen may be a gas which consists of H.sub.2 and N.sub.2. Pure H.sub.2 (100% H.sub.2) may also be used. The diffused hydrogen terminates the surface states and reduces junction leakage.\nFIG. 1 is a sectional view of a portion of a conventional semiconductor device which includes a transistor 10 formed on a semiconductor substrate 1. Transistor 10 may, for example, be formed in a peripheral circuit region of a semiconductor memory device such as a dynamic random access memory device (DRAM) or may be part of a logic circuit. Transistor 10 includes a source/drain diffusion region 2 and a drain/source diffusion region 3 which are spaced apart by a channel region 4. A gate electrode 6 is insulatively spaced from channel region 4 by a gate insulating layer 5. Source/drain region 2 and drain/source region 3 may be formed, for example, by performing an ion implantation using gate electrode 6 as an implantation mask. An isolation region 7 of an insulator such as silicon dioxide (SiO.sub.2) electrically isolates transistor 10 from other elements formed on semiconductor substrate 1.\nIn high density large scale integrated circuit (LSI) devices such as 64 Mbit and 256 Mbit DRAMs, a self-aligned contact (SAC) process is often utilized for forming device contacts. Using an SAC process, the spacing between the contacts can be reduced and a high integration density can be achieved. In many SAC processes, a silicon nitride layer is formed over the LSI device area to serve as an etch-stop layer. As a result, with reference to FIG. 1, a silicon nitride layer 8 is formed over isolation region 7, source/drain region 2, drain/source region 3, and gate electrode 6. However, silicon nitride layer 8 is substantially a diffusion barrier and thus prevents the diffusion of hydrogen into semiconductor substrate 1. Thus, in high density LSIs in which a silicon nitride layer is formed during an SAC process, the diffusion of hydrogen into semiconductor substrate 1 to reduce junction leakage is prevented. This is a particular problem, for example, for transistors in the peripheral circuit region of a DRAM or for transistors in a logic circuit. While hydrogen could be diffused into the substrate before the silicon nitride layer is formed, the hydrogen can be easily removed from the surface states by high temperature LSI process steps which follow the formation of the silicon nitride layer. Thus, the hydrogen diffusion is most effectively carried out after the high temperature steps of the LSI fabrication process. However, as described above, at this time, the silicon nitride layer acts as a barrier to such diffusion."}
{"text": "This invention relates to an ignition timing regulating device for an internal combustion engine, which senses the pressure of ambient atmosphere at, for example, the engine starting time and regulates the ignition timing of the engine on the basis of the value of sensed atmospheric pressure.\nIt is necessary to determine the ignition timing of an internal combustion engine depending on the state of the engine so that the engine can operate in its optimum operation mode. As is commonly known, it is the best mode of engine operation to ignite the engine under the condition of minimum advance at maximum torque or so-called MBT (minimum advance for best torque) when the factors such as the efficiency and fuel consumption of the engine are taken into account. To this end, it is necessary to change or regulate the ignition timing depending on the state of the engine so that the ignition can always occur under the condition which satisfies the MBT. This MBT is a function of the conditions of the engine and is primarily determined by main and auxiliary factors. The rotation speed of the engine and the pressure in the intake manifold are representative of the main factors, and the temperature of engine cooling water, EGR, etc. are representative of the auxiliary factors. The value of MBT is also variable depending on the value of atmospheric pressure. The value of MBT varies depending on the value of atmospheric pressure for the reason that, when the engine operates on, for example, a highland, the value of atmospheric pressure is lower than that on a lowland, that is, the volume density of air is reduced, with the result that the air-fuel ratio (A/F) in the carburetor shifts toward a richer side of the setting. According to the results of a test on a four-cylinder four-cycle internal combustion engine, the value of MBT varies relative to the value of A/F at the engine rotational frequency of 2000 rpm and the intake manifold pressure of -200 mm Hg in a manner as shown in FIG. 1. It is therefore necessary to regulate the advance so that it meets the variation in the value of MBT.\nAn engine ignition timing control device of electronic type has been proposed hitherto to regulate the advance so that it meets the value of MBT which varies in a manner as shown in FIG. 1. In the proposed ignition timing control device, the length of time required for the engine to make one complete revolution or to rotate through a predetermined angle is measured to detect the rotation speed of the engine, and a semiconductor type pressure sensor or a pressure responsive diaphragm is used to electrically sense the intake manifold pressure. In such an ignition timing control device, the value of advance .alpha..sub.N based on the detected rotation speed of the engine N and the value of advance .alpha..sub.P based on the sensed intake manifold pressure P are added to provide the sum (.alpha..sub.N +.alpha..sub.P), and on the basis of this sum, the ignition timing is primarily determined. Then, various connection factors are added to finally determine the ignition timing. An atmospheric pressure sensor of diaphragm type, semiconductor type or any other suitable type is additionally employed to sense the pressure of ambient atmosphere so as to retard the advance angle of the ignition timing depending on the value of sensed atmospheric pressure. This atmospheric pressure sensor is basically the same in function as the intake manifold pressure sensor in that both of them are pressure sensors. Hitherto, both of these two pressure sensors have been provided from the standpoint that they serve individually the separate functions.\nThe prior art ignition timing control device of the type above described is however expensive in that it includes the atmospheric pressure sensor and the intake manifold pressure sensor in spite of the fact that these two sensors are basically the same in their functions. Further, because of the recent demands for purification of engine exhaust gases, low fuel consumption, etc., it is required to sense the atmospheric pressure with high accuracy and to suitably compensate the value of sensed intake manifold pressure for the control of the ignition timing. However, a relatively inexpensive atmospheric pressure sensor comprising the combination of a diaphragm and a switch is capable of only sensing the pressure at a single point from the structural aspect, and a plurality of such sensors corresponding to the number of pressure sensing points are required to sense the pressures at such points. Further, employment of a pressure sensor of semiconductor or like type in the device for the purpose of improvement in the accuracy of pressure sensing leads inevitably to a further higher cost. Thus, an inexpensive device which can sense the atmospheric pressure with high accuracy is now demanded."}
{"text": "The business of a call center, also known as a customer center, is to provide rapid and efficient interaction between agents and customers (or prospective customers). Existing solutions require the purchase of multiple hardware and software components, typically from different vendors, to achieve the business goals of the customer center. The use of separate systems of components leads to a variety of problems. For instance, each system typically has its own method of configuration and its own user interface. Thus, exchanging data between the systems requires additional work by someone at the customer center.\nFurthermore, customer centers are continually tasked with striking a balance between service quality, efficiency, effectiveness, revenue generation, cost cutting, and profitability. As a result, today's customer center agents are charged with mastering multiple data sources and systems, delivering consistent service across customer touch points, up-selling, cross-selling, and saving at-risk customers, while winning new ones.\nAdditionally, more and more layers of complexity resulting from regulations and standards add further to the challenges of operating a customer center. For example, federal regulations like Sarbanes-Oxley or HIPAA, or business rules such as product return policies, are becoming increasingly challenging to track and monitor, particularly the effectiveness agents' in meeting compliance objectives.\nThus, there remains a need for an automation technology to drive, for example and among others, adherence, compliance, alerts, triggers, and workflows."}
{"text": "1. Field of the Invention\nThe present invention relates to flash EPROM memory technology, and more particularly to unique cell structure having extended floating gates to improve the coupling ratio between the control gate, floating gate, and source or drain of the transistor cell.\n2. Description of Related Art\nFlash EPROMs are a growing class of non-volatile storage integrated circuits. These flash EPROMs have the capability of electrically erasing, programming, and reading a memory cell in the chip. The memory cell in a flash EPROM is formed using so-called floating gate transistors, in which the data is stored in a cell by charging or discharging the floating gate. The floating gate is a conductive material, typically made of polysilicon, which is insulated from the channel of the transistor by a thin layer of oxide, or other insulating material, and insulated from the control gate or word line of the transistor by a second layer of insulating material.\nData is stored in the memory cell by charging or discharging the floating gate. The floating gate is charged through a Fowler-Nordheim tunneling mechanism by establishing a large, positive voltage between the control gate and source or drain. Alternatively, an avalanche mechanism may be used by applying potentials to induce high energy electrons in the channel of the cell which are injected across the insulator to the floating gate. The voltage on the control gate is divided by the so-called coupling ratio of the cell, resulting in a first voltage between the control gate and floating gate, and a second voltage between the floating gate and the source or drain. With a 50% coupling ratio, half of the voltage applied to the control gate appears across the oxide between the floating gate and the source or drain. This voltage between the floating gate and the source or drain causes electrons to tunnel or be injected into the floating gate through the thin insulator. When the floating gate is charged, the threshold voltage for causing the memory cell to conduct is increased above the voltage applied to the word line during a read operation. Thus, when a charged cell is addressed during a read operation, the cell does not conduct. The non-conducting state of the cell can be interpreted as a binary 1 or a zero, depending on the polarity of the sensing circuit.\nThe floating gate is discharged to establish the opposite memory state. This function is typically carried out by F-N tunneling between the floating gate and the source or drain of the transistor, or between the floating gate and the substrate. For instance, the floating gate may be discharged through the source by establishing a large, positive voltage from the source to the gate, while the drain is left at a floating potential.\nThe high voltages used to charge and discharge a floating gate place significant design restrictions on flash memory devices, particularly as the cell dimensions and process specifications are reduced in size. Thus, the coupling ratio for the memory cells becomes a critical design parameter.\nOne way of increasing the coupling ratio, is to increase the surface area of the floating gate between the control gate and the floating gate. This can be accomplished by extending the floating gate over the source or drain regions, such as described in Bergemont, et al., U.S. Pat. No. 5,012,446.\nOne approach to extending the floating gate is described in Kume, et al., \"A 1.28 .mu.m.sup.2 Contactless Memory Cell Technology for a 3 V-only 64 Mbit EEPROM\", IEDM 92, pp. 991-993; Hisamune, et al., \"A High Capacitive-Coupling Ratio (HiCR) Cell for 3 V-only 64 Mbit and Future Flash Memories\", IEDM 93, pp. 19-22.\nAll of these prior art designs for extending the floating are relatively complicated process technologies.\nAccordingly, an improved process for extending the floating gate to increase the coupling ratio of flash EPROM, and a circuit for utilizing such structure, is desirable."}
{"text": "Traditionally, typical network computing environments include point-to-point communication between two computing devices. However, some communications technologies employ multi-point communications, where groups of computing devices simultaneously receive a common transmission.\nUnicast and Multicast\nThere are at least two common approaches to simultaneously transmitting the same content to multiple computing devices over a network computing environment: unicast or multicast.\nUnicast may be understood to be a communication that takes place over a network between a single sender and a single receiver. With unicast, a computing device generates and sends one set of data packets (each has common content) to each receiving computer. As the receiving group increases in size, unicast becomes increasingly less efficient because it is simultaneously transmitting copies of the same data packets to an increasing number of computing devices. In direct and geometric proportions, unicast requires increasingly more bandwidth as the receiving group increases. That is because the same information is carried multiple times—even on shared links.\nMulticast may be understood to be a communication that transmits a single message to a select group of multiple recipients. In contrast to broadcasting, multicasting typically refers to sending a single message to a select group on a network rather than to everyone connected to the network.\nWith multicast, a computing device generates only one copy and sends it to the select group that chooses to receive it. This technique addresses packets to a group of receivers rather than to a single receiver. It typically depends on the network infrastructure to forward the packets to only the sub-networks and the receivers that need to receive them.\nA common implementation of multicasting is Internet Protocol (IP) Multicast. It is a bandwidth-conserving technology that reduces traffic by simultaneously delivering a single stream (e.g., a media stream) of information to thousands of recipients. Typical applications that take advantage of multicast include videoconferencing, corporate communications, distance learning, and distribution of software, stock quotes, and news.\nMedia Streams\nWith the advent of digital media streaming technology (such as those using IP multicast), users are able to see and hear digital media, more or less, as the data is being received from a media server.\nHerein, a “media stream” is a multimedia object (containing audio and/or visual content, such as a video) that is compressed and encoded in accordance with mechanisms generally available now or in the future for doing so. Furthermore, such a media stream is intended to be decoded and rendered in accordance with generally available mechanisms for doing so.\nWithout a loss of generality, the same techniques can be applied to any media stream that has a similar structure which reduces temporal, spatial, or perceptual redundancies. For example, many audio compression formats such as AC3 have keyframes followed by modification data to regenerate an approximation of the original uncompressed stream.\nMultimedia Distribution Format Standards\nDue to the amount of data required to accurately represent such multimedia content, it is typically delivered to the computing device in an encoded, compressed form. To reproduce the original content for presentation, the multimedia content is typically decompressed and decoded before it is presented.\nA number of multimedia standards have been developed that define the format and meaning of encoded multimedia content for purposes of distribution. Organizations such as the Moving Picture Experts Group (MPEG) under the auspices of the International Standards Organization (ISO) and International Electrotechnical Commission (IEC), and the Video Coding Experts Group (VCEG) under the auspices of the International Telecommunications Union (ITU) have developed a number of multimedia coding standards (e.g., MPEG-1, MPEG-2, MPEG-4, H.261, H.263, and the like).\nThere are many different standardized video-stream data formats. For example: MPEG, H.263, MPEG-1, MPEG-2, MPEG-4 Visual, H.264/AVC, and DV formats. Likewise, there are many different standardized audio-stream data formats. For example: MPEG audio, AC3 audio, DTS audio, or MLP audio.\nMPEG-2/H.262\nThe predominant digital video compression and transmission formats are from a family called block-based motion-compensated hybrid video coders, as typified by the ISO/IEC MPEG-x (Moving Picture Experts Group) and ITU-T VCEG H.26x (Video Coding Experts Group) standards. This family of standards is used for coding audio-visual information (e.g., movies, video, music, and such) in a digital compressed format.\nFor the convenience of explanation, the MPEG-2 video stream (also known as an H.262 video stream) is generally discussed and described herein, as it has a structure that is typical of conventional video coding approaches. However, those who are skilled in the art understand and appreciate that other such digital media compression and transmission formats exist and may be used.\nAn example representation of a MPEG-2 format is shown in FIG. 1. Each video sequence is composed of a sequence of frames that is typically called Groups of Pictures (or “GOP”), such as GOP 105. A GOP is composed of a sequence of pictures or frames. The GOP data is compressed as a sequence of I-, P- and B-frames where:                An I-frame (i.e., intra-frame) is an independent starting image—(compressed in a similar format to a JPEG image). An I-frame or “key frame” (such as I-frame 110t) is encoded as a single image, with no reference to any past or future frames. It is sometimes called a random access point (RAP).        A P-frame (i.e., forward predicted frame) is computed by moving around rectangles (called macroblocks) from the previous I- or P-frame then (if so indicated by the encoder) applying a ‘correction’ called a residual. Subsequent P-frames (such as P-frame 120t) are encoded relative to the past reference frame (such as a previous I- or P-frame).        Zero or more B-frames (i.e., bi-directional predicted frames, such as frames 130 and 132) are formed by a combination of rectangles from the adjacent I- or P-frames, followed (if so indicated by the encoder) by a correction residual.        \nThe GOP structure is intended to assist random access into the stream. A GOP is typically an independently decodable unit that may be of any size as long as it begins with an I-frame.\nTransmission and Presentation Timelines\nFIG. 1 illustrates two manifestations of the same MPEG-2 video stream. The first is the transmission timeline 100t and the other is the presentation timeline 100p. This is an example of transmission and presentation timelines of a typical media stream and their relationship to each other.\nThe transmission timeline 100t illustrates a media stream from the perspective of its transmission by a media-stream encoder and transmitter. Alternatively, it may be viewed from the perspective of the receiver of the transmission of the media stream.\nAs shown in FIG. 1, the I-frame (e.g., 110t) is typically temporally longer than the other frames in the transmission timeline. Since it doesn't utilize data from any other frame, it contains all of the data necessary to produce one complete image for presentation. Consequently, an I-frame includes more data than any of the other frames. Since the I-frame has more data than others, it follows that it typically requires greater time for transmission (and, of course, reception) than the other frame types.\nFIG. 1 also shows P-frames (such as 120t) and B-frames (such as 130t and 132t) of the transmission timeline 100t. Relative to the B-frames, the P-frames are temporally longer in the transmission timeline because they typically include more data than the B-frames. However, P-frames are temporally shorter than I-frames because they include less data than I-frames. Since the B-frames rely on data from at least two other frames, they typically do not need as much data of their own to decode their image as do P-frames (which rely on one other frame).\nFIG. 1 also illustrates the presentation timeline 100p of the media stream from the perspective of its presentation by the media decoder and presenter. In contrast to their transmission duration, the presentation duration of each frame—regardless of type—is exactly the same. In other words, it displays at a fixed frequency.\nThe incoming frames of the media stream are decoded, buffered, and then presented at a fixed frequency (e.g., 24 frames per second (fps)) to produce a relatively smooth motion picture presentation to the user. In MPEG 2 used to convey NTSC video, the field rate is fixed, and each MPEG 2 picture may produce 1, 2, or 3 fields. Field pictures are required to produce 1 field, and frame pictures may produce 2 or 3 fields. Thus, the frame picture presentation rate may not be fixed, but it is not dictated by the transmission rate of the frame pictures.\nFIG. 1 also illustrates a typical decoded GOP 105 of MPEG in its presentation timeline. This GOP example includes an I-frame 110p; six P-frames (e.g., 120p); and 14 B-frames (e.g., 130p and 132p). Typically, each GOP includes a series of consecutively presented decoded frames that begin with an I-frame (such as frame 110p).\nGOP Presentation Delay\nFIG. 1 shows that the I-frame 110t of an example GOP is first received beginning at point T1 in time; however, it is not first presented until point T2. The time gap between the two points is called herein the “GOP presentation delay” and is labeled 170 in FIG. 1. It represents the delay from when the receiver first begins receiving the first frame of a GOP (which is typically the I-frame) until the device first presents the first frame of the GOP.\nMedia-Stream Presentation Start-up Delay\nTo tune channels in a media-streaming environment (such as in a multicast environment), a receiver requests a target channel. It receives the target media stream and then waits for an access point into the stream. A channel change cannot occur until an access point is received. From the perspective of the user, this can lead to lengthy channel change times.\nFIG. 2 illustrates an example of a media-stream presentation start-up delay at 280. The start-up delay is the effective delay experienced by a user. It includes a delay between when a particular media stream is requested and when the first frame of a GOP from the particular media stream is actually presented. As shown in FIG. 2, the start-up delay 280 includes the GOP presentation delay 270 (discussed above).\nReferring to FIG. 2, this example is explained. A GOP, starting with I-frame 210t, is being transmitted. This is shown in the transmission timeline 200t. The receiver seeks to tune into this media stream at request point R. This selection is illustrated as a user selecting a media-stream channel using a remote control 260.\nAgain, this is an example illustration for explanatory purpose. This point R could be at any moment in time after the beginning (i.e., after the beginning of its I-frame 210t) of a GOP.\nThe receiver must wait for a random access point (or RAP) in order to access the media stream. In this example, each GOP has one RAP. An I-frame is an example of a typical RAP. Therefore, each GOP has one I-frame. So, the receiver must wait for the next I-frame (at the beginning of the next GOP) before it can access the media-stream transmission as shown by transmission timeline 200t. \nOnce the receiver has an I-frame in its buffer, it may refer back to it for dependency decoding of P- and B-frames. Consequently, a conventional system must wait for a RAP before it can start buffering frames (that are useful).\nIn FIG. 2, the receiver starts buffering the next GOP at point M1 with I-frame 250t. Thus, the first frame that may be eventually presented to the user is I-frame 250t, because it is the first RAP in the stream after the point at which the receiver joined the stream. Because of the GOP presentation delay (discussed above), it actually starts presenting the GOP (with I-frame 250p of presentation timeline 200p) at point M2—which is also the presentation start-up point S of the start-up delay 280.\nAs demonstrated by the screens 262-266, the start-up delay is the effective delay experienced by a user. The user selects a media-stream channel at request point R (using, for example, a remote 260) and sees a blank screen, as shown by screen 262. Of course, there may be information presented here (such as electronic programming information), but, since it is not yet the desired media-stream content, it is effectively blank.\nScreen 264 shows that screen remains blank even after the next GOP is currently being received. Screen 266 shows that the first image of frame 250p is finally presented to the user.\nThe average length of this start-up delay is directly affected by the average GOP length. Some media-stream providers employ relatively long average GOP lengths. In those instances, this delay is even more acute because, when changing channels, the user is waiting longer for the next GOP to come around.\nIt short, this start-up delay is very annoying to, and provokes impatience in, the typical users."}
{"text": "It is a common occurrence in telephone conversations for one party to place another party on hold while the former party performs tasks, such as referencing records, answering another telephone call and the like. Many telecommunication systems play music during the duration of the holding period to provide a pleasurable distraction for the holding party.\nAlthough many parties prefer to listen to music during the holding period, some parties would prefer to suppress or deactivate the music, particularly where three or more parties are engaged in a conference call. In conference calls involving three or more parties, the music can impede the holding parties' ability to engage in a conversation during the holding period.\nExisting telecommunication infrastructures do not have the capability of suppressing or deactivating the music, unless the source of the music is local to the telecommunication equipment serving the holding party. It is often not possible for the remote telecommunication equipment to process deactivation commands from the local switch."}
{"text": "The present invention relates to gas turbine cycles employing steam injection for power augmentation. More specifically, a system and method are provided that are capable of regulating steam production as required by modern gas turbines with Dry Low NOx (DLN) combustion systems.\nSteam injection for gas turbine power augmentation, including steam generation with exhaust heat recovery has been known for many decades. See, e.g., U.S. Pat. No. 2,678,531. The basic concept shown in this patent is widely used even now with non-DLN combustion systems where a once-thru steam generator is used for steam production with exhaust heat. Steam injection at part load gas turbine operation is allowed in non-DLN systems, which make such a once-thru steam generator system a good design for such applications. However, DLN combustion systems impose several additional restrictions on the operating regimes in which steam injection is allowed. These new DLN requirements are that steam injection is permitted only when the gas turbine is at base load and that steam injection is only permitted above a certain ambient temperature, typically above 59xc2x0 F. The steam generation system should thus have the flexibility to operate with no steam production for all gas turbine operations up to base load, operate at base load with no steam production (when the steam injection is not permitted due to low ambient temperature), and be able to regulate steam production from minimum turndown steam production rate to the maximum steam demand at base load when steam injection is permitted. Meeting these requirements adds considerable capital costs, operating costs and complexity to the exhaust heat recovery steam generation systems known currently. once-thru steam generation system has limited turndown capability (approximately 3:1), and reducing the flow below this minimum limit can result in flow instability, vibration, and tube failure. When the turbine operating regime does not allow for steam injection (in DLN systems), the steam would have to be vented to the atmosphere. Since high quality demineralized water is used in a once-thru steam generator, venting the steam when steam injection is not allowed would result in considerable increase in operating cost, and would be an unattractive option to end users of such power plants.\nThere are additional potential design options that may be considered. For example, completely turning-off the water flow/steam production in the once-thru boiler, when there is no requirement for steam injection, has been considered. The design of a once-thru steam generator for dry operation (i.e., with no flow inside the tubes, while hot gases flow outside) at maximum gas turbine exhaust temperatures (xe2x88x921200xc2x0 F.) requires the use of high cost alloy materials. Also, the introduction of cold water into the evaporator tubes of a dry hot boiler (when there is steam demand following dry operation) results in thermal shock which would result in considerable reduction in the life of the boiler. The reliability and cost impact make such a design unattractive.\nOne might also consider condensing the steam and re-use of the water when there is no demand for steam. However, such steam condensing heat exchangers would have to be manufactured with alloy material compatible with the high steam temperature resulting in high capital cost for the addition of such equipment.\nThe addition of a steam turbine to expand unused steam has been disclosed in U.S. Pat. No. 5,727,377. Condensing the expanded low temperature steam and re-use of the water is also discussed in this patent. Additional equipment required in such a system, additional capital cost, and added complexity in control make such a system unattractive for peaking power application.\nDiverting the hot exhaust gas to a parallel duct installed with the heat recovery section has been disclosed in the U.S. Pat. No. 3,693,347, to control the steam production rate in a once-thru boiler. This would require a diverter damper. Modern gas turbines have very high exhaust gas flow rates and such diverter dampers become extremely large, unreliable, and costly. Also, large diverter dampers always tend to have leakage and such leakage could result in the stagnant gas volume in the heat transfer section approaching the exhaust gas temperature over long term operation. Hence, the heat transfer section material would have to be selected similar to the dry operation described above, and there is also the potential for thermal shock of the hot tubes when water is introduced.\nThe use of drum type boilers for exhaust heat recovery and varying the drum pressure to regulate the steam production rate has also been disclosed by others (U.S. Pat. Nos. 5,566,542 and 4,393,649). While steam production rate can be regulated by varying the drum pressure, high exhaust temperatures of modern gas turbines would result in a very limited turndown capability from such a control method.\nThus, the DLN restrictions have resulted in the need to develop an exhaust heat recovery steam generation system with improved steam production regulation capability, for the economical implementation of steam injection in modern gas turbine power plants with DLN combustion systems.\nThe invention is embodied in a steam generation system and method that provides an economical solution for implementation of exhaust heat recovery steam generation for injection to modern gas turbines with DLN combustion systems. More specifically, the invention is embodied in a unique and improved method of controlling the steam production rate that has the capability to control the steam production rates from no steam generation to maximum rated steam production at base load gas turbine operation and no steam production at all, part load gas turbine operation. The invention provides this flexibility while being able to use conventional low cost material for heat recovery sections and a heat rejection system.\nThus, the invention is embodied in a gas turbine cycle that comprises a gas turbine system including a compressor for compressing air, a combustion system for receiving compressed air from the compressor and a turbine for converting the energy of the combustion mixture that leaves the combustion system into work; a heat recovery system for receiving exhaust gas from the gas turbine; a first flow path for at least one of water and steam at elevated pressure including a first heat exchange flow path disposed in heat exchange relation to the exhaust gas flowing through the heat recovery system, thereby to produce an at least partly evaporated fluid stream; a second flow path for at least one of water and steam defining a power augmenting flow path operatively coupled to the first flow path for receiving flow therefrom, the second flow path including a second heat exchange flow path disposed in heat exchange relation to the exhaust gas flowing through the heat recovery system, thereby to produce superheated steam, the second flow path being operatively coupled to at least one of the gas turbine compressor discharge and the combustion system of the gas turbine for increasing the mass flow of fluid thereinto; a third flow path for at least one of water and steam defining a recirculating flow path operatively coupled to the first flow path for receiving fluid therefrom and for recirculating the fluid to the first heat exchange flow path; and a heat rejector system for selectively cooling fluid flowing through the third flow path thereby to reduce a temperature of the fluid upstream of the first heat exchange flow path.\nIn an exemplary embodiment, a separator, more specifically a dearator, is provided for separating the at least partly evaporated fluid stream into saturated steam for flow through the second flow path and water for recirculation through the third flow path.\nIn one embodiment, the heat rejection system comprises parallel flow paths defined by a split of the third flow path, and a heat exchanger is disposed along one of the parallel flow paths for reducing the temperature of the fluid flowing therethrough. A valve unit may be disposed at an upstream juncture of the parallel flow paths for controlling flow of fluid into each of the parallel flow paths.\nThe second heat exchange flow path is advantageously disposed downstream, in a direction of exhaust gas flow, from the first heat exchange flow path.\nThe invention is also embodied in a method for augmenting the power produced by a gas turbine system of the type having a compressor for compressing air to produce compressed air, a combustor for heating said compressed air and producing hot gases, and a turbine for receiving said hot gases and converting the energy thereof to work for driving said compressor, for supplying a load, and for producing hot exhaust gases, the method comprising flowing a fluid through a first flow path including a first heat exchange flow path disposed in heat exchange relation to said exhaust gas, thereby to produce an at least partly evaporated fluid stream; selectively flowing a part of said fluid stream through a second flow path including a second heat exchange flow path disposed in heat exchange relation to said exhaust gas, thereby to produce superheated steam for injection to at least one of the gas turbine compressor discharge and the combustion system of the gas turbine for increasing the mass flow of fluid thereinto; selectively flowing a remainder of said fluid stream through a third flow path for recirculating said fluid to said first heat exchange flow path; and selectively cooling fluid flowing through said third flow path thereby to reduce a temperature of said fluid upstream of said first heat exchange flow path."}
{"text": "Multi-processor systems include many types of processors. The processors of a server are often housed together in a single structure, creating high-density systems with a modular architecture that ensures flexibility and scalability. Thus, this modular architecture reduces space requirements. Server processors, along with storage, networking and other types of processors, are typically installed in a common enclosure, or chassis, that hosts multiple processors that share common resources such as cabling, power supplies, and cooling fans.\nMulti-processor systems create challenging engineering problems, due largely in part to heat produced by the processors and limited space in the chassis. Typically, multi-processor systems are limited by an underlying power and thermal envelope. For example, a chassis that hosts a multi-processor system may only be designed to utilize a limited number of watts. That is, the chassis can only consume so much power and is limited in the amount of airflow that is available to cool the processors in the chassis.\nEngineering challenges occur in optimizing the tradeoff between performance and thermal and power requirements. In a multi-processor system multiple processors, each representing a separate system, are present in the same chassis. Associated with the chassis are a specific set of power and thermal requirements. Specifically, these requirements put a limit on the amount of power that can be consumed by the respective processors. Known power limiting strategies include powering down a CPU functional unit, e.g., a floating-point unit or an on-die cache, or trading off speed for reduced power consumption in a hard drive.\nThis power limitation puts a constraint on the frequency that the processors can run, and thus, limits the performance. In addition, the processors in a system are usually all configured to operate at the same frequency. This further limits the ability for the individual processors to operate at optimal performance and capacity.\nPrior solutions run all the processors at a performance level less than their maximum in order to meet the overall chassis power and thermal cooling budget. A disadvantage associated with this solution is that the performance of each processor is degraded or diminished to fall within these budgets. For example, if the ability of the chassis to cool is limited to X and there are Y processors, each processor can only contribute approximately X/Y to the dissipated power in the chassis. Thus, each processor is limited to the performance associated with an X/Y power level.\nAnother solution has been to add a plurality of loud, space-consuming fans that require expensive control circuitry. These cooling systems increase the cost of the multi-processor system, leave less space for other features within the chassis for other features, and run a higher risk for failures and increased downtime. Other solutions have included limiting the number of I/O cards in the multi-processor system, as well as restricting the number of other use features. A further solution has been to reduce the power budget available for other features in the multi-processor system.\nWhat is needed is a method for optimizing the performance of a multi-processor system by modulating the voltage of the processors within the system in conjunction with the frequency of the processors to increase the thermal and power benefit of the multi-processor system."}
{"text": "1Field of the Invention\nThe present invention relates generally to electrical switches of the snap-action type and, more particularly to a novel snap switch which operates satisfactorily at very low current and voltage levels and which includes a contact which has a self-cleaning action.\n2. Background Art\nSnap switches are well known devices for switching the flow of electricity and are manufactured, for example by Westport Development & Manufacturing Company, Inc. The type of such switches to which the present invention relates have a one-piece spring blade with a contact carrying free end which is selectively moved between electrical terminals. The blade is of the snap-acting, over center type and has an operation point which, if sufficiently linearly moved beyond a prescribed limit, distorts the blade to effect the snap-action and a subsequent movement of the free contact end from one terminal to another. The movement of the blade is maintained within a range by two spaced apart stop members. The amount of movement of the operation point from one stop member until the blade shifts from one contact to the other is precisely determined and is generally maintained constant throughout a wide range of operating temperatures.\nSuch type of snap switches is described in U.S. Pat. Nos. 3,278,700; 3,609,269; and 4,017,699, issued to Robert R. Hellman, and assigned to the assignee of the present invention, the disclosures of all of which patents are hereby made a part hereof by reference. The first of the patents describes a snap switch having a bracket for preloading the blade and certain features for temperature compensation. The second of the patents describes such a switch disposed in a hermetically sealed housing and actuatable by means of the distortion of a sidewall of a tube forming part of the housing. The third of the patents describes a snap switch employing two snap switches mounted together so that both are actuated simultaneously.\nWhile the switches described above have provided satisfactory service, they are somewhat limited in that there are lower limits of voltage and current below which they cannot be used, due to fouling of the flat surfaces contacted by the round contact ball at the free end of the blade. This fouling is due, in part, to the creation of polymer from ambient impurities by the contact ball as it only slightly wipes the surfaces. This slight wiping action is insufficient to clean the polymer thus formed from the surfaces. In addition, when the contact ball contacts a spherical surface on a terminal, there is a flattening of the contacting surfaces because of the wiping action. The lower limits of known such switches are about 10 milliamperes and 10 volts.\nAnother factor contributing to the limitations of the prior snap switches is that the blade itself is in the electrical circuit between two terminals, since the fixed end of the blade is mounted on a common terminal. This introduces additional resistance into the total circuit, since desirable spring materials are relatively poor electrical conductors. The lowest total switch resistance obtained with conventional switches has been on the order of about 60-80 milli-Ohms.\nFurthermore, it is desirable that the current carrying components of such switches have the most efficient electrical characteristics, but that these be combined with lower cost materials for other components requiring other characteristics such as weldability and suitability for high temperature service.\nAccordingly, it is an object of the present invention to provide a snap switch that is capable of use at lower voltage and current levels than previously known such switches.\nAn additional object of the invention is to provide such a switch that employs a self-cleaning contact.\nAnother object of the invention is to provide such a switch in which the blade thereof does not become part of the electrical circuit being switched.\nA further object of the invention is to provide such a switch which is usable at higher temperatures than previously known such switches.\nYet an additional object of the invention is to provide such a switch that may be economically constructed.\nYet another object of the invention is to provide such a switch which has lower total electrical resistance than previously known such switches.\nOther objects of the present invention, as well as particular features and advantages thereof, will be elucidated in, or be apparent from, the following description and the accompanying drawing figures."}
{"text": "1. Field of the Invention\nThe invention relates to apparatus for providing control over both a transmission of a vehicle and an implement, such as a snow plow, mounted on that vehicle.\n2. Description of the Prior Art\nOftentimes, a plow, particularly a snow plow, is mounted on the front of a general purpose vehicle, such as a truck or tractor, and various driver actuated controls for raising, lowering and/or angling the plow are appropriately situated in the cab of the vehicle. To effectively plow an area, the driver faces the tasks of controlling both the movement of the vehicle and the position of the plow. Usually, to control the vehicle, the driver must maintain one hand on its steering wheel and the other on its transmission lever. Generally, the driver actuated plow controls, which typically comprise one or more levers, are usually mounted below or underneath the dashboard of the vehicle. Typically, to re-position the plow, the driver must first stop the vehicle, in order to take his hand off the transmission lever, and then lean down and appropriately operate the snow plow controls. Thereafter, the driver returns to his normal seated position to set the vehicle in motion. If a relatively large unobstructed area, such as a highway, is being plowed, then the tasks of controlling both the vehicle and the plow are relatively effortless. The driver firsts positions the plow and then starts the vehicle moving. In such a case, relatively little time is spent positioning the plow and the majority is spent controlling the movement of the vehicle. By contrast, these tasks become very time-consuming and fatiguing whenever relatively small or densely packed areas, such as residental streets or driveways, are to be plowed. Here, to avoid obstructions, the driver must constantly stop the vehicle, re-position the plow, and then restart the vehicle in motion. Unfortunately, in this situation, the driver wastes considerable time and is easily fatigued through excess repetitive efforts which entail stopping the vehicle, taking his hand off the transmission lever, bending down, operating the plow controls, returning to his normal driving position and then re-starting the vehicle in motion by shifting gears, steering and appropriately depressing the accelerator. It will be understood that with some prior art controls it is not always necessary to stop the vehicle to re-position the plow. Nonetheless, the driver must still take one hand off the gear shift or transmission lever to re-position the plow, and, if a change in gears is required while the plow is being re-positioned, this can not be readily accomplished with these prior art controls.\nVarious arrangements exist in the art to control the position of a plow. See, for example, that disclosed in U.S. Pat. Nos. 4,221,266 (issued to Fardal on Sept. 9, 1980); 4,026,048 (issued to Hill et al on May 31, 1977) and 3,585,319 (issued to Payerle on June 15, 1971). Other arrangements, typified by that disclosed in U.S. Pat. No. 3,941,009 (issued to Brown on Mar. 2, 1976), exist for shift levers to control vehicular transmissions. However, none of these arrangements minimizes the fatigue associated with operating a snow plow by eliminating the excess repetitive efforts required of a driver."}
{"text": "1. Field of the Invention\nThe present invention relates to a microprocessor that can prevent illegal alternation of execution codes and processing target data under a multi-task program execution environment.\n2. Description of the Background Art\nIn recent years, the performance of a microprocessor has improved considerably such that the microprocessor is capable of realizing reproduction and editing of video images and audio sounds, in addition to the conventional functions such as computations and graphics. By implementing such a microprocessor in a system designed for end-user (which will be referred to as PC hereafter), the users can enjoy various video images and audio sounds on monitors. Also, by combing the function for reproducing video images and audio sounds with the computational power of the PC, the applicability to games or the like can be improved. Such a microprocessor is not designed for any specific hardware and can be implemented in a variety of hardwares so that there is an advantage that the users who already possess PCs can enjoy reproduction and editing of video images and audio sounds inexpensively by simply changing a microprocessor for executing programs.\nIn the case of handling video images and audio sounds on PCs, there arises a problem of a protection of the copyright of original images or music. In the MD or digital video playback devices, unlimited copies can be prevented by implementing a mechanism for preventing the illegal copying in these devices in advance. It is rather rare to attempt the illegal copying by disassembling and altering these devices, and even if such devices are made, there is a worldwide trend for prohibiting the manufacturing and sales of devices altered for the purpose of illegal copying by laws. Consequently, damages due to the hardware based illegal copying are not very serious.\nHowever, image data and music data are actually processed on the PC by the software rather than the hardware, and the end-user can freely alter the software on the PC. Namely, if the user has some level of knowledge, it is quite feasible to carry out the illegal copying by analyzing programs and rewriting the executable software. In addition, there is a problem that the software for illegal copying so produced can be spread very quickly through media such as networks, unlike the hardware.\nIn order to resolve these problems, conventionally a PC software to be used for reproducing copyright protected contents such as commercial films or music has employed a technique for preventing analysis and alternation by encrypting the software. This technique is known as a tamper resistant software (see David Aucsmith et al., “Tamper Resistant Software: An Implementation”, Proceedings of the 1996 Intel Software Developer's Conference).\nThe tamper resistant software technique is also effective in preventing illegal copying of valuable information including not only video and audio data but also text and know-how that is to be provided to a user through the PC, and protecting know-how contained in the PC software itself from analysis.\nHowever, the tamper resistant software technique is a technique which makes analysis using tools such as de-assembler or debugger difficult by encrypting a portion of the program that requires protection before the execution of the program starts, decrypting that portion immediately before executing that portion and encrypting that portion again immediately after the execution of that portion is completed. Consequently, as along as the program is executable by a processor, it is always possible to analyze the program by carrying out the analysis step by step starting from the start of the program.\nThis fact has been an obstacle for a copyright owner to provide copyright protected contents to a system for reproducing video and audio data using the PC.\nThe other tamper resistant software applications are also vulnerable in this regard, and this fact has been an obstacle to a sophisticated information server through the PC and an application of a program containing know-how of an enterprise or individual to the PC.\nThese are problems that equally apply to the software protection in general, but in addition, the PC is an open platform so that there is also a problem of an attack by altering the operating system (OS) which is intended to be a basis of the system's software configuration. Namely, a skilled and malicious user can alter the OS of his own PC to invalidate or analyze the copyright protection mechanisms incorporated in application programs by utilizing privileges given to the OS.\nThe current OS realizes the management of resources under the control of the computer and the arbitration of their uses by utilizing a privileged operation function with respect to a memory and an execution control function provided in CPU. Targets of the management include the conventional targets such as devices, CPU and memory resources, as well as QoS (Quality of Service) at network or application level. Nevertheless, the basics of the resource management are still allocations of resources necessary for the execution of a program. Namely, an allocation of a CPU time to the execution of that program and an allocation of a memory space necessary for the execution are the besics of the resource management. The control of the other devices, network and application QoS is realized by controlling the execution of a program that makes accesses to these resources (by allocating a CPU time and a memory space).\nThe OS has privileges for carrying out the CPU time allocation and the memory space allocation. Namely, the OS has a privilege for interrupting and restarting an application program at arbitrary timing and a privilege to move a content of a memory space allocated to an application program to a memory of a different hierarchical level at arbitrary timing, in order to carry out the CPU time allocation. The latter privilege is also used for the purpose of providing a flat memory space to the application by concealing (normally) hierarchical memory systems with different access speeds and capacities from the application.\nUsing these two privileges, the OS can interrupt an execution state of the application and take a snap shot of it at arbitrary timing, and restart it after making a copy of it or rewriting it. This function can also be used as a tool for analyzing secrets hidden in the application.\nIn order to prevent an analysis of the application on a computer, there are several known techniques for encrypting programs or data (Hampson, U.S. Pat. No. 4,847,902; Hartman, U.S. Pat. No. 5,224,166; Davis, U.S. Pat. No. 5,806,706; Takahashi et al., U.S. Pat. No. 5,825,878; Buer et al., U.S. Pat. No. 6,003,117; Japanese Patent Application Laid Open No. 11-282667 (1999), for example). However, these known techniques do not account for the protection of the program operation and the data secrecy from the above described privileged operations of the OS.\nThe conventional technique based on the x86 architecture of Intel Corporation (Hartman, U.S. Pat. No. 5,224,166) is a technique for storing the execution codes and data by encrypting them by using a prescribed encryption key Kx. The encryption key Kx is given in a form of EKr[Kx] which is encrypted by using a public key Kp corresponding to a secret key Ks embedded in a processor. Consequently, only the processor that knows Ks can decrypt the encrypted execution codes on a memory. The encryption key Kx is stored in a register inside the processor called a segment register.\nUsing this mechanism, it is possible to protect the secrecy of the program codes from the user to some extent by encrypting the codes. Also, it becomes cryptographically difficult for a person who does not know the encryption key Kx of the codes to alter the codes according to his intention or newly produce codes that are executable when decrypted by using the encryption key Kx.\nHowever, the system employing this technique has a drawback in that the analysis of the program becomes possible by utilizing a privilege of the OS called a context switching, without decrypting the encrypted execution codes.\nMore specifically, when the execution of the program is stopped by the interruption or when the program voluntarily calls up a software interruption command due to the system call up, the OS carries out the context switching for the purpose of the execution of the other program. The context switching is an operation to store an execution state (which will be referred to as a context information hereafter) of the program indicating a set of register values at that point into a memory, and restoring the context information of another program stored in the memory in advance into the registers.\nFIG. 15 shows the conventional context storing format used in the x86 processor. All the contents of the registers used by the application are contained here. The context information of the interrupted program is restored into the registers when the program is restarted. The context switching is an indispensable function in order to operate a plurality of programs in parallel. In the conventional technique, the OS can read the register values at a time of the context switching, so that it is possible to guess most of the operations made by the programs if not all, according to how the execution state of that program has changed.\nIn addition, by controlling a timing at which the exception occurs by setting of a timer or the like, it is possible to carry out this processing at arbitrary execution point of the program. Apart from the interruption of the execution and the analysis, it is also possible to rewrite the register information by malicious intention. The rewriting of the registers can not only change the operation of the program but also make the program analysis easier. The OS can store arbitrary state of the application so that it is possible to analyze the operation of the program by rewriting the register values and operating the program repeatedly. In addition to the above described functions, the processor has a debugging support function such as a stepwise execution, and there has been a problem that the OS can analyze the application by utilizing all these functions.\nAs far as data are concerned, U.S. Pat. No. 5,224,166 asserts that the program can access the encrypted data only by the program execution using the encrypted code segment. Here, there is a problem that the encrypted data can be freely read by the encrypted program by using arbitrary key, regardless of the encryption key by which the program is encrypted, even when there are programs encrypted by using mutually different encryption keys. This conventional technique does not account for the case where the OS and the application have their own secrets independently and the secret of the application is to be protected from the OS or a plurality of program providers have their own secrets separately.\nOf course, it is possible to separate memory spaces among the applications and to prohibit accesses to a system memory by the applications by the protection function provided in the virtual memory mechanism even in the existing processor. However, as long as the virtual memory mechanism is under the management of the OS, the protection of the secret of the application cannot rely on the function under the management of the OS. This is because the OS can access data by ignoring the protection mechanism, and this privilege is indispensable in providing the virtual memory function as described above.\nAs another conventional technique, Japanese Patent Application Laid Open No. 11-282667 (1999) discloses a technique of a secret memory provided inside the CPU in order to store the secret information of the application. In this technique, a prescribed reference value is required in order to access data in the secret memory. However, this reference fails to disclose how to protect the reference value for obtaining the access right with respect to the secret data from a plurality of programs operating in the same CPU, especially the OS.\nAlso, in U.S. Pat. No. 5,123,045, Ostrovsky et al. disclose a system that presupposes the use of sub-processors having unique secret keys corresponding to the applications, in which the operation of the program cannot be guessed from the access pattern by which these sub-processors are accessing programs placed on a main memory. This is based on a mechanism for carrying out random memory accesses by converting the instruction system for carrying out operations with respect to the memory into another instruction system different from that.\nHowever, this technique requires different sub-processors for different applications so that it requires a high cost, and the implementation and fast realization of the compiler and processor hardware for processing such instruction system are expected to be very difficult as they are quite different from those of the currently used processors. Besides that, in this type of processor, it becomes difficult to comprehend correspondences among the data contents and the operations even when the data and the operations of the actually operated codes are observed and traced so that the debugging of the program becomes very difficult, and therefore this technique has many practical problems, compared with the other conventional techniques described above in which the program codes and the data are simply encrypted, such as those of U.S. Pat. No. 5,224,166 and Japanese Patent Application Laid Open No. 11-282667."}
{"text": "The present invention pertains to optical communications. More particularly, this invention relates to an optical multi-plexing device with passively aligned molded optics which spatially disperses collimated multi-wavelength light from an optical fiber into individual wavelength bands. In certain preferred embodiments, the improved multiplexing device of the present invention is particularly suited for wavelength division multiplexing systems for the fiber optic data-communications and telecommunications systems.\nIn wavelength division multiplexed optical communication systems, many different optical wavelength carriers provide independent communication channels in a single optical fiber. Future computation and communication systems place ever-increasing demands upon communication link bandwidth. It is generally known that optical fibers offer much higher bandwidth than conventional coaxial communications; furthermore, a single optical channel in a fiber waveguide uses a microscopically small fraction of the available bandwidth of the fiber (typically a few GHz out of several tens of THz). By transmitting several channels at different optical wavelengths into a fiber [i.e., wavelength division multiplexing (WDM)], this bandwidth may be more efficiently utilized.\nPrior art optical multiplexers and demultiplexers include the Nosu et al U.S. Pat. No. 4,244,045. Nosu utilizes a series of wavelength sensitive filters arranged in a zigzag optical pathway at a predetermined oblique angle to the light beam. Nosu positions filters on both sides of substrate 60. It is significant that the fabrication of the Nosu device is expensive and time consuming. As described at column 4, lines 33-52, the multi-layer filters are separately grown on the substrate and have as many as ten layers each. The couplers and lenses must be separately aligned in a tedious, expensive post-fabrication process.\nThe Scobey U.S. Pat. Nos. 5,583,683 and 5,786,915 teach an eight channel multiplexing device in which a continuous variable thickness interference filter is deposited onto the surface of an optical block. This design has inherent weaknesses. First, each filter must necessarily integrate the signal over its width, since its thickness (and wavelength sensitivity) varies across its width; resulting in less precise filtering. Secondly, the interference filter thickness may vary from device to device, also degrading its filtering performance. Thirdly, each of the individual couplers such as item 62 of FIG. 3 must be separately aligned after fabrication.\nThe Scobey et al U.S. Pat. No. 5,859,717 teaches the use of a precision optical block which is opaque (col. 4, lines 20-29) and in which slots or gaps are machined or formed to provide an optical path. Any machining required adds to the expense of producing the recesses, gaps and slots. Furthermore, the separate collimators (6,24,46) and filters 32 must be separately aligned after fabrication of the device, adding to the time and expense of production.\nThe Jayaraman U.S. Pat. No. 5,835,517 teaches an optical multiplexer having multiple optical cavities which must each be \"tuned\" after fabrication by adjusting its length. Such \"tuning\" adds significantly to cost of production, and is expressly avoided with the present invention.\nThe Lemoff et al U.S. Pat. 5,894,535 teaches an optical multiplexer in which a plurality of dielectric channel waveguides are embedded in a cladding region, which process is considerably more expensive than the monolithic formed parts of the present invention. Furthermore, the filters must be separately aligned with the channel waveguides after fabrication. Perhaps the greatest inherent disadvantage of the Lemoff et al design is the difficulty, expense and time of mounting the mirror 36 in channel 50 and properly mounting the filters 45a through 45d. The channel 50 and the filter mounting surface must be cut with an expensive excimer laser or microtome, yielding a relatively rough mounting surface. The rough mounting surface creates bonding and alignment problems. The filters, for example, are angularly sensitive and must be mounted precisely vertically and parallel with each other as shown in FIG. 3. Additionally, since the waveguides are so broad, the waveguides have large angular content, which causes a \"rolled off\" or variable filter response. The present invention avoids the bonding and alignment problems of Lemoff et al as well as the variable filter response. The present invention uses preformed and/or premolded extremely flat surfaces for carrying the reflective mirror surface and the filters, so that the mirror and filters are properly bonded and aligned. The present invention has much smaller divergence which, together with proper filter alignment, results in much cleaner separation of the n wavelength beams by the filters.\nAn optical multiplexing device combines or separates multiple light signals with varying optical frequencies. The optical multiplexing device has applications for both dense and coarse wavelength division multiplexing (DWDM and CWDM) for both multi-mode and signal-mode fiber optic data communications and telecommunications. Multiple wavelength light sources are combined into a single optical path for transmission or multi-wavelength light travelling in a single optical path is separated into multiple narrow spectral bands that is focused onto individual fiber optic carriers or detectors.\nCurrent wavelength division multiplexed (WDM) devices are designed for operation in single-mode optical fiber telecommunications systems, where performance over long distances (100 km) is the primary factor and cost and size are secondary. As bandwidth demands within the Networking Industry [Local Area Networks (LAN) and Wide Area Networks (WAN)] increase, compact inexpensive wavelength division multiplexed systems will become necessary in order to utilize the full bandwidth of the optical fiber. The WDM device described herein utilizes plastic-mold injection and inexpensive dielectric filters to create a compact device capable of multiplexing or demultiplexing multiple optical wavelengths.\nOne of the key features and a primary object of the present invention is to provide a compact and cost effective optical multiplexer and demultiplexer for both single-mode or multi-mode fiber optic communication systems wherein the device includes preformed and premolded passively aligned optics.\nA further object of the invention is to provide an optical multiplexer/demultiplexer having only two preformed parts and a series of filters, the filters being sandwiched between the preformed parts.\nAnother feature of the present invention is to minimize optical loss due to divergence of light between the source coupling optics and the fiber optic connector coupling optics.\nAnother feature of the present invention is to integrate a fiber optic connector within the optical module to passively align an optical fiber to the coupling optic.\nA further feature of the present invention is to allow demultiplexed multi-wavelength light to couple directly to a photodetector array without any intervening optical fiber.\nThe present invention is comprised of a single plastic molded coupling module, a single precision optical block with a reflective coating on one side, an array of discrete multi-wavelength Fabry-Perot transmission filters in one embodiment, and index matching optical adhesive to bond said components together.\nThe plastic molded coupling module is formed by integrating an aspherical off-axis collimating lens array, filter mechanical alignment features, a redirectional mirror, a coupling lens, a fiber optic connector, and mechanical features for passive alignment of the optical part to either an array of light sources, a detector array, or an optical fiber array, all within a single part.\nA precision optical block can be formed of any transparent optical material capable of transmitting light over the desired spectral region and being formed or polished to the correct thickness. A reflective surface can be placed on the optical block by a number of techniques, to include; dielectric interference coatings, metallic coatings, etc."}
{"text": "In the past, it has been particularly difficult to provide a completely waterproof submersible lamp which is able to withstand knocks and vibrations in use while being able to maintain waterproof integrity. This problem is particularly manifest in the tail lights of boat trailers but can also be present in other motor vehicle situations, and in other locations such as in underwater lights in swimming pools or in areas which are often subjected to high pressure water treatment.\nIn a trailer-like situation, the problem is made worse by the necessity to provide different colored lenses within a single light. For example, the tail and brake lights of a motor vehicle must have a red lens, the turn indicator light must have an amber lens and there is also often the necessity to provide a clear lens section for illuminating the adjacent number plate of the trailer or other motor vehicle. It has proven extremely difficult to manufacture conventional combined tail/stop/indicator/number plate lights which will maintain waterproof integrity at the joints of the various different colored lamp segments. It is of course possible to provide individual lamps for each of the above functions but this significantly increases the cost and also the complexity and \"untidiness\" on the motor vehicle.\nThe legal requirements of many countries also make it necessary for the various lenses to have certain light focusing or diffracting characteristics and it has proven extremely difficult to provide these characteristics in a waterproof tail light or indicator light lens of the type described above.\nThere are also many situations where it is desirable to provide a waterproof entry for an electrical cable into an apparatus. In the past, in some instances, an elastomeric collar or O-ring is provided which forms an interference fit with the outer sheath of the cable to inhibit the influx of water. In other cases, the cable passes through a bush or collar and is \"potted\" in place using a setting chemical compound such as a silicon based compound or an epoxy resin. Such waterproof entries can work loose over time, particularly in situations subject to high vibration or stress loadings which will then allow the ingress of water into the apparatus. Furthermore, if the outer sheath of a multi-conductor electrical cable becomes ruptured or pierced, water can enter into the housing through the cable sheath, even if the entry point remains intact about the external perimeter of the sheath."}
{"text": "Marine and land-based industrial (M & I) gas turbine engines are frequently derived from engines designed for and used in various types of aircraft. Such marine and industrial gas turbine engines are used, for example, for powering marine vessels, electrical generators, and various types of pumping and compression applications.\nThe parent gas turbine engine of an M & I engine is typically designed and constructed to be lightweight and to operate at minimum specific fuel consumption (SFC) in an aircraft for predetermined thermodynamic cycles of operation having predetermined ranges of air and combustion gas flowrates, temperature, and pressure in the engine. The parent, or aircraft, engine cycles are also preselected for maximizing thermal and propulsive efficiency of the engine.\nDevelopment of an aircraft gas turbine engine requires a substantial amount of design, development, and testing resulting in substantial development costs. In designing gas turbine engines for marine and industrial applications, it has proved to be more cost-effective to modify an existing aircraft gas turbine engine in the desired power class, than to design the M & I engine from the beginning. Accordingly, it is desirable to minimize the changes in the aircraft engine required for obtaining a suitable M & I engine.\nOne application of an M & I engine is to provide peaking power for powering an electrical generator to provide additional electrical power to a utility power grid when the utility power demands exceed on-line baseload capacity. Accordingly, the utility industry desires relatively simple and inexpensive gas turbine engines which can be brought on line quickly and then shut down quickly as required to meet the peaking requirement. One goal is to generate the required power during peaking operation as efficiently as possible for reducing kilowatt-hour costs.\nOne factor in obtaining relatively low kilowatt-hour cost is the development costs for providing an industrial gas turbine engine for meeting the required power demands. In order to keep development costs relatively low, the industrial gas turbine engine typically utilizes the parent aircraft gas turbine engine and makes as few changes in the design thereof as practical for obtaining the desired land-based gas turbine engine. Accordingly, the parent aircraft gas turbine engine utilized for M & I applications may be adapted specifically for particular applications, including, for example, driving an electrical generator at a synchronous speed such as 3000 rpm or 3600 rpm, for generating electricity at 50 Hertz or 60 Hertz, respectively. Since electrical power generation is desired, and in particular, for meeting peaking utility demands, maximum output shaft horsepower from the gas turbine engine is desirable for maximizing the amount of electrical power generated by the generator."}
{"text": "Telephones have typically received the small amount of power required for their operation directly from the telephone lines. With such telephone sets, the telephone may be operable even though the primary line power to an establishment is interrupted.\nMore recently, microprocessor based electronics have been included in telephone sets to provide for telephone functions which are based on digital processing. For example, a central processing unit (CPU) based on a microprocessor may provide for automatic dialing from a directory of electronically stored numbers or from temporary storage of the last number dialed by the telephone set. Where the telephone set is to be used for transmitting both voice and digital data signals, the CPU may control routing of the signals to and from speakers and microphones and to and from digital processing circuitry. Where the telephone set services two or more telephone lines, the CPU may control selection of a single line or of both lines for conference calls.\nThe power available on the telephone line is insufficient for the digital processing circuitry of the more sophisticated telephone sets. Thus, the electronics rely on the normal local power. To prevent noise on the phone line from the power line, the circuits powered by the phone line and local power line must be isolated from each other through transformers, optical couplers and the like. The isolation is most effective and simplest when made at or close to the incoming telephone lines. As a result, substantially all of the telephone circuit is powered by the local power line. A disadvantage of such systems is that the telephone set becomes dependent on both telephone lines and power lines for operation. Simpler telephones may be operable with interruptions in power, but the sophisticated electronic telephones become inoperable."}
{"text": "The instant invention provides methyl substituted-2-oxohexane derivatives and processes for preparing same as well as processes for using same for augmenting or enhancing the aroma of perfume compositions, colognes and perfumed articles. The methyl substituted-2-oxohexane derivatives of our invention are defined according to the structure: ##STR3## wherein R.sub.12 represents hydroxyl having the structure: EQU --OH]\nor acetoxy having the structure: ##STR4## and wherein R.sub.3, R.sub.4, R.sub.5, R.sub.6 and R.sub.7 each represent the same or different methyl or hydrogen with the provisos that:\n(i) the sum total of the carbon atoms of R.sub.3, R.sub.4, R.sub.5, R.sub.6 and R.sub.7 is three; PA1 (ii) when R.sub.7 is methyl, then R.sub.5 and R.sub.6 are each methyl; and PA1 (iii) when R.sub.7 is hydrogen, then R.sub.3 or R.sub.4 is methyl. PA1 (i)--Murphy & Lane, Ind. Eng. Chem., Prod. Res. Dev., Vol. 14, No. 3, 1973, p. 167 (Title: Oligomerization of 2-Methyl-2-Butene in Sulfuric Acid and Sulfuric-Phosphoric Acid Mixtures). PA1 (ii)--Whitmore & Mosher, Vol. 68, J. Am. Chem. Soc., February, 1946, p. 281 (Title: The Depolymerization of 3,4,5,5-Tetramethyl-2-Hexene and 3,5,5-Trimethyl-2-Heptene in Relation to the Dimerization of Isoamylenes). PA1 (iii)--Whitmore & Stahly, Vol. 67, J. Am. Chem. Soc., December, 1945, p. 2158 (Title: The Polymerization of Olefins. VIII The Depolymerization of Olefins in Relation to Intramolecular Rearrangements. II). PA1 (iv)--U.S. Pat. No. 3,627,700, issued on Dec. 14, 1971, (Zuech). PA1 (v)--U.S. Pat. No. 3,538,181, issued on Nov. 3, 1970, (Banks). PA1 (vi)--U.S. Pat. No. 3,461,184 issued on Aug. 12, 1969 (Hay, et al). PA1 (vii)--Tausz, Chemische Zentralblatt, 1919, II, 125, (Production of Di-isoamylene From Isoamylene Using Mercury Acetate Catalyst).\nChemical compounds which can provide woody, camphoraceous, patchouli-like, fruity and amber-like aromas are highly desirable in the art of perfumery. Many of the natural materials which provide such fragrances and contribute such desired nuances to perfumery compositions and perfumed articles are high in cost, unobtainable at times, vary in quality from one batch to another and/or are generally subject to the usual variations in natural products.\nThere is, accordingly, a continuing effort to find synthetic materials which will replace, enhance or augment the fragrance notes provided by natural essential oils or compositions thereof. Unfortunately, many of the synthetic materials either have the desired nuances only to a relatively small degree or they contribute undesirable or unwanted odor to the compositions.\nApplication for U.S. Letters Patent Ser. No. 184,132 filed on Sept. 4, 1980 now U.S. Pat. No. 4,321,255 issued on Mar. 23, 1982 (incorporated by reference herein) discloses the production of unsaturated branched ketones according to the reaction: ##STR5## wherein in each of the structures containing dashed lines, these structures represent mixtures of molecules wherein in each of the molecules, one of the dashed lines represents a carbon-carbon double bond and each of the other of the dashed lines represent a carbon-carbon single bond. These compounds so produced are indicated to be useful for their organoleptic properties in augmenting or enhancing the aroma or taste of consumable articles including perfume compositions and colognes. It is noteworthy that the compounds defined according to the structure: ##STR6## have unsaturation in the structure and, in addition, contain twelve carbon atoms rather than ten carbon atoms. Disclosure of the resulting alcohols produced by reaction of these ketones with methyl magnesium halide is set forth in application for U.S. Letters Patent Ser. No. 212,887 filed on Dec. 4, 1980 (now U.S. Pat. No. 4,318,934 issued on Mar. 9, 1982). The resulting alcohols have thirteen carbon atoms rather than the ten carbon atom containing alcohols of the instant application.\nNothing in the prior art discloses compounds defined according to the generic structure: ##STR7## wherein R.sub.1, R.sub.2, R.sub.3, R.sub.4, R.sub.5, R.sub.6 and R.sub.7 are defined supra.\n\"Diisoamylene\" is indicated to be synthesized in the following references:\nUnited Kingdom Pat. No. 796,130 published on June 4, 1958 discloses the synthesis of polyalkylindanes by means of, interalia, reacting alpha-methylstyrene with trimethylethene (2-methyl-butene-2) in the presence of an acid catalyst such as sulfuric acid or boron trifluoride methyletherate. It is further indicated that such compounds are useful intermediates in the production of perfumery compounds. Apparently, however, the more volatile diisoamylenes produced as side-products in the reaction of 2-methyl-butene-2 with alpha-methylstyrene are discarded.\nThe diisoamylene useful as a starting material in the instant case may be distilled from the reaction product (as see Example A infra) at a temperature in the range of 36.degree.-40.degree. C.; a liquid temperature in the range of 74.degree.-94.degree. C. and a pressure of 4-5 mm/Hg."}
{"text": "Fixation of a toner on a print is achieved, among other things, by rolling a heated fixing roller with pressure on the print and reliably joining the toner in this way to the print. In order to achieve reliable separation of the fixing roller from the print, a fixing oil subsequently also referred to in general as a fixing aid, is often applied to the fixing roller. A shortcoming is that the fixing oil is taken up in an undesired manner by the printer, and covers parts of the printer. Because of this, the printer result is also adversely affected, if the fixing oil reaches the printed image. Cleaning passes of the printer, to free it of the fixing aid, are therefore prescribed in the prior art. However, it is unspecified when and whether a cleaning pass is necessary."}
{"text": "1. Field of Invention\nThis invention relates to the technical field of mimeograph, and more particularly to a method of controlling a linkage drive section in a stencil printing machine to set a printing drum in place in which the position of a print image can be adjusted in a horizontal direction with ease. In this specification, the term \"horizontal direction\" is intended to mean \"a direction which is in parallel with the central axis of the printing drum\"; and the term \"horizontal position\" is intended to mean \"a position along the direction thus defined\".\n2. Description of Related Art\nIn a conventional rotary stencil printing machine with a printing drum, the printing drum is set at a predetermined position in the direction of its axis in the printing machine frame and rotated around its central axis. In the machine, adjustment of the horizontal position of a print image can be achieved by suitably shifting the position of the mimeographic stencil paper on the printing drum in the direction of the axis of the printing drum before it is wound on the printing drum. Alternatively, by moving the sheet supplying device in the direction of the axis of the printing drum, the horizontal position of the print image can be adjusted at all times, not only in the case where the stencil paper is not wound on the printing drum yet but also in the case where it is wound on it.\nIn moving a mimeographic stencil paper or a printing sheet, which is rectangular and has two short sides extended in the direction of movement and two long sides which are perpendicular to the short sides, over the printing drum along the axis of the latter, the stencil paper or printing sheet must be parallel-moved in the direction of the short sides with high accuracy. In order to make this movement delicate and accurate thereby to permit the fine adjustment in horizontal position of the print image, means for horizontally moving a stencil paper feeding device adapted to feed a stencil paper to the printing drum and a sheet supplying device adapted to supply printing sheets to the printing drum should have: guide means which use a pair of parallel guide rails laid in the lateral direction of the stencil paper or printing sheet to guide the stencil paper feeding device or the sheet supplying device accurately parallel at least two positions spaced from each other in the longitudinal direction of the stencil paper or printing sheet; and fine drive means for giving feeding actions to the stencil paper feeding device or the sheet supplying device along the guide means for the same distance at the same speed.\nAs is apparent from the above description, in the case where, with respect to the printing drum which has been fixed in the axis direction, the stencil paper or printing sheet is moved along the axis of the printing drum to adjust the horizontal position of the print image, it is necessary to provide considerably intricate means."}
{"text": "The present invention relates to a fixing device."}
{"text": "In keeping up with the recent tendency towards using multimedia equipment, a need has been felt to process a larger quantity of the information, including the video information, and to increase the capacity of a storage device used for recording the information. In particular, in the field of storage of the high-quality video information, the tendency is towards using a BD (Blu-ray Disc) which is larger in storage capacity than the DVD (Digital Versatile Disc). To increase the recording capacity of an optical disc or a HDD device, it is necessary to increase the recording density. Concomitantly, reducing an error rate to provide for sufficient reliability has become a crucial task. The following description is centered mainly about increasing the density of the optical disc.\nIn an optical disc, a light beam collected by an optical element is illuminated on a disc medium, and the information is detected by lightness/darkness of reflected light or by light polarization. A beam spot collected is finite and, the smaller its diameter, the higher may be the density in recording/reproduction. In this situation, an optical approach towards reducing the beam spot is progressing. The spot diameter is inversely proportionate to the NA (Numerical Aperture) of an objective lens and directly proportionate to the wavelength λ of the laser beam. The spot diameter may thus be reduced by increasing the NA and by decreasing the wavelength λ. However, if NA is increased, the depth of focus becomes shallower, such that it becomes necessary to decrease the distance between the disc surface and the lens. Thus, limitations are imposed on increasing the NA. On the other hand, in a short wavelength laser, stability in high output oscillation and long useful life are requirements. Although the wavelength used is becoming shorter, as evidenced by use of an infrared laser (λ=780 nm) for a CD, a red laser (λ=650 nm) for DVD and a blue laser (λ=405 nm) for BD, limitations are again imposed on further reducing the wavelength.\nMeanwhile, the MTF (Modulation Transfer Function), which is the frequency characteristic of a transmission route between an optical head and a disc medium, is in the form of an LPF (Low-Pass Filter) in which the gain in the high frequency range is decreased by reason of the finite beam spot. Hence, even if a rectangular wave is recorded, the readout waveform from a disc becomes dull. If the recording density is increased, a waveform to be read out at a specified time point interferes with a wavelength to be read out at another time point, a phenomenon known as inter-symbol interference. By reason of this inter-symbol interference, a recorded mark shorter than a preset length becomes difficult to reproduce. If conversely the recording mark is long, the frequency with which the phase information used for extracting the clock for synchronization is output is decreased, thus causing pulling-out of synchronization. It is therefore necessary to limit the mark length to less than a certain length. The recording data for the optical disc is encoded for recording from the above described perspective of signal processing. In particular, the RLL code (Run Length Limited Code) in which limitations are imposed on the length between transitions is preferentially used. For example, EFM (Eight to Fourteen Modulation), an 8/16 code or (1, 7) PP is used. The minimum run length of the PFM modulation code used for a CD and that of the 8/16 modulation code used for a DVD are 2 (d=2). The minimum run length of the (1, 7) PP, used in a BD of higher density, is 1. This code is 2/3 as is the (1, 7) RLL, and is featured by limitations imposed on the number of succession of shortest marks.\nThere is also known a technique termed waveform equalization. This technique uses an inverse filter that eliminates the inter-symbol interference. The technique emphasizes the high frequency components of the readout signal and hence suppresses the inter-symbol interference. However, the high frequency components of the noise are simultaneously emphasized, thus possibly deteriorating the SNR (Signal to Noise Ratio). In particular, if the recording density is increased, the deterioration of the SNR caused by waveform equalization is mainly responsible for errors in data detection. The PR (Partial Response) equalization is a system of waveform equalization that intentionally causes known inter-symbol interference to occur on purpose. In this technique, high frequency components are usually not emphasized, thus suppressing the SNR from deterioration.\nAnother effective detection system is the maximum likelihood detection system, according to which detection is carried out on a data sequence known to undergo certain state transitions. The detection performance may be improved by selecting such a time-series pattern, out of all possible time-series patterns, which has a smallest value of a mean square sum of errors. However, if the above processing is carried out on real circuits, difficulties are encountered in connection with the circuit size and the operational speed. Hence, an algorithm termed a Viterbi algorithm is usually employed to progressively select the path. In this case, the Viterbi algorithm is termed the Viterbi decoding or Viterbi detection.\nA detection system, which is a combination of the above mentioned PR equalization and Viterbi detection, is known as PRML (Partial Response Maximum Likelihood) system, and detects data as it performs a sort of error correction. By PR equalization, correlation along the time direction is imparted to the readout signal. Hence, only specified state transitions are presented in the data sequence obtained on sampling the readout signal. The state transitions, thus specified, and a data sequence of the noise-corrupted actual readout signal, are compared to each other, and most probable state transitions are selected, whereby errors in data detection may be reduced. A PRML detection system, in which the modulation code with the minimum run length and a PR (1, 2, 2, 2, 1) channel are used, is described in Non-Patent Document 1. With this detection system, a broader detection margin may be obtained in high density recording/reproduction.\nTo improve the detection performance by Viterbi detection, it is necessary that the frequency characteristic of the read channel is made to be coincident with a specified PR equalization characteristic. In this case, such PR equalization characteristic that is as close to that of the read channel as possible is to be selected. In general, the frequency characteristic is corrected using a waveform equalizer so that the PR characteristic will be as close to the preset PR characteristic as possible. Among the techniques of adaptively correcting signal deterioration with time to improve the detection performance, there are an automatic equalization system and an adaptive equalization system. Among the algorithms of adaptive equalization of the sequential type, there are a zero forcing method and a mean square method. The adaptive equalization technique is highly advantageous in that initial device adjustment is unneeded.\nThe operation of Viterbi detection will now be described. FIG. 12 shows signal transition states of a signal xn obtained by sampling a readout signal of a DVD in timed relation to the channel clock followed by PR (1, 3, 3, 1) equalization. FIG. 13 depicts a trellis diagram in which state transitions of FIG. 12 are plotted along the time axis. The numbers affixed to branches extending from six states of the state transition diagram represent ideal amplitude values rn of xa. For example, in an area composed of 4T spaces and 4T marks in succession in this order, the state sequence is S0→S1→S3→S7→S6→S4→S0→S0→S1→ . . . . xn at this time is ideally −4, −3, 0, 3, 4, 3, 0, −3, −4, −3, . . . . However, due to e.g., the noise, xa may become −3.9, −2.9, 0.1, 2.7, 3.8, 2.9, 0.2, −2.6, −3.9, −3.1, . . . . Now suppose that the Gaussian noise be superposed on xn and rn, assume any one of five reference levels (±14, ±3, 0), it is then the maximum likelihood detection to find rn that minimizes Σ(xn−rn)2. It is however difficult to compare all combinations in real time. In this consideration, the Viterbi algorithm sequentially performs the operation. As shown in FIG. 12, possible previous states of S0 are S0 and S4, and possible previous states of S7 are S3 and S7. Then referring to FIG. 14, the operation of selecting more plausible one of a plurality of paths entered simultaneously, such as states S0 and S7, is carried out at each time point. By tracing back the paths towards the past, the paths merge at a certain time point in one path (path merge). In short, the information prior to the merge time point may be identified.\nFor selecting one of the paths, an index for plausibility, termed a metric, is introduced. The plausibility Pan that a time point n is a state Sa is termed a path metric. A square of the difference between xn and the reference level r is termed a branch metric bn(r).bn(r)=(xn−r)2  (1)\nSince the path metric is an integration of the branch metrics from the past, the smaller the path metric, the more plausible is the path metric. Since the path metric P1n is necessarily the state S0 one time point before, the path metric P1n is a path metric P0n-1, which is state S0 one time point before, plus the branch metric bn (−3) at the present time point, as indicated by the equation (2):P1n=P0n-1+bn(−3)  (2)\nIn similar manner, P3n, P4n and P6n are as indicated by the following equations (3) to (5):P3n=P1n-1+bn(0)  (3)P4n=P6n-1+bn(0)  (4)P6n=P7n-1+bn(3)  (5)\nP0n is indicated by the equation (6). P0n is S0 or S4 one time point before, such that there are two paths. Out of the path metrics, a smaller one is selected. It should be noted that Min[a, b] denotes a or b, whichever is smaller.P0n=Min[P0n-1+bn(−4),P4n-1+bn(−3)]  (6)\nIn similar manner, P7n becomes as indicated by the following equation (7):P7n=Min[P7n-1+bn(4),P3n-1+bn(3)]  (7)\nThe path metrics are updated at every time point to select the path. The paths entered in all states merge in one path. Hence, by tracing back from a given time point towards the past, the paths merge to determine the information. Meanwhile, in the equations 6 and 7, it is sufficient that large-small comparison may be made between the path metrics. Since the terms xn2 in bn(r) are common to all path metrics, the following equation:bn(r)′=r2−2rxn may be used as a branch metrics, whereby the circuit may be simplified to advantage.Meanwhile, this processing is usually carried out in terms of a channel clock as a unit. Hence, a high speed processing is required. For example, at a speed eight times as fast as the DVD, the channel clock frequency exceeds 200 MHz. Hence, the processing is generally carried out by a dedicated circuit. FIG. 15 depicts a block diagram of a Viterbi detector. A path metric-branch metric addition circuit (Add), a path metric value comparator circuit (Compare) and a path selection circuit that selects the path based on the result of comparison (Select) are collectively termed an ACS circuit 52. Further, a circuit for calculating the branch metric 51 (BMG) and a memory (MEM) 53 that holds the path selection information, termed a path memory, are needed. It has been known that, if the signal level distribution after adaptive equalization for each reference level is proximate to the normal distribution, PRML detection by a Viterbi detector is able to detect even a low SNR signal that may not be detected by level detection. The low SNR signal not being detected by level detection may be exemplified by such a case where a signal has a good recording mark quality but suffers from significant jitter, a case where a signal suffers from significant inter-symbol interference, and a case where a signal has been reproduced in the defocused state.\nThere are however cases where the performance of a PRML detection system is inferior to that of the conventional level detector. These cases may be met when a mark non-optimum in shape has been recorded on a disc mainly due to failure (unmatching) of power control in recording. FIG. 6 shows a typical readout waveform of a 5T space+11T mark, in which the mark is a long mark suffering from distortion. Such mark has been formed by a recording strategy in which the recording power is to be reduced at the leading end of the long mark. In the level detection, an edge (transition) position of 5T space to −11T mark is closer to the 5T side, and hence may be detected without errors by bi-level slice detection. On the other hand, suppose now that the edge position be detected by PRML. In this case, one out of an ideal path (5T+11T, correct solution) and an ideal path 2 (6T+10T, error) is selected. Thus, path metrics of the two paths are now to be found. In Table 1, square values of the differences between the ideal values and x, at each time point, viz., branch metric values, are shown for the ideal paths 1, 2 in the PR (1, 3, 3, 1) channel.\nTABLE 1ideal path 1ideal path 2time nxn(rn)(xn − rn)2(r2n)(xn − r2n)23−4.1−40.01−40.014−3.5−30.25−40.255−1.101.21−33.6161.532.2502.2572.0443182.243.2443.24Σ——10.96—10.36\nThe sum of the branch metrics gives a path metric. In a Viterbi detector, one of the two paths having a smaller path metric value is determined to be more plausible. The path metric for the ideal path 1 is 10.96, whereas that for the ideal path 2 is 10.36. Hence, the path 2 with the smaller path metric is selected, which results in an error detection.\nFIG. 8 shows a typical readout waveform of a 6T space+10T mark whose long mark suffers from distortion. Such mark has been formed by a recording strategy in which the recording power is to be increased at the leading end portion of the long mark. In the level detection, an edge position of 6T space to 10T mark is closer to the 10T mark side, and hence may be detected without errors by bi-level slice detection. On the other hand, suppose now that this be detected by PRML. In this case, an ideal path (5T+11T, error) or an ideal path 2 (6T+10T, correct solution) is to be selected. Path metrics of the two paths are now to be found. In Table 2, square values of the differences between the ideal values and x, at each time point, viz., branch metric values, are shown for the ideal paths 1, 2 in the PR (1, 3, 3, 1) channel.\nTABLE 2ideal path 1ideal path 2time nxn(rn)(xn − rn)2(r2n)(xn − r2n)23−4.1−40.01−40.014−3.5−30.25−40.255−1.201.44−33.2461.034.001.074.540.2532.2584.140.0140.01Σ——5.96—6.76\nThe sum of the branch metrics gives a path metric. In the Viterbi detector, one of two paths which has a smaller path metric value is determined to be more plausible. The path metric for the ideal path 1 is 5.96, while that for the ideal path 2 is 6.76. Hence, the path 1 with the smaller path metric is selected, which results in an error detection.\nThe level detection is a conventional technique used in many optical disc drives. Even in an optical disc drive, provided with a Viterbi detector, which has come to be used these days, reproduction compatibility with respect to a disc recorded with a conventional drive is indispensable. It is true that, in detection by PRML, the reproduction performance may be expected to be improved. However, in detecting a readout signal with strong non-linearity, such as a case with a non-optimum recording mark, the detection performance may become inferior to that with threshold detection, which is counted as a drawback.\nPatent Document 1 discloses a measure to be taken in avoiding the deterioration in the performance in the detection by PRML of a signal suffering from waveform distortion. FIG. 16 shows an arrangement of the Viterbi detector. An ACS circuit 102 takes charge of addition of path metrics and branch metrics generated by a branch metric generation circuit 101, comparison of path metric values, and path selection which is based on the results of comparison. A path memory 103 processes the path selection information. A maximum likelihood decision unit 104 selects and outputs data which represents the smallest path metric value. However, upon path selection, the information obtained on comparing an input signal entered to the Viterbi detector at a time point a preset number of channel clocks before, and a threshold value α, to each other by a comparator 106 is used to make the latest path decision, in addition to the path metric values. In particular, if the input signal at the time point a preset number of channel clocks before has an amplitude corresponding to that of a mark, such a path where the marks are consecutively arrayed is selected. If otherwise, path selection is made on the basis of the path metric values. By such path selection, such a path where marks are consecutively arrayed may be selected even in case there is mark distortion, thus prohibiting an error in data detection.\nPatent Document 2 discloses an alternative method for detecting the waveform distortion by PRML. FIG. 17 depicts a block diagram showing its formulation. In this method, the information recorded on a magneto-optical recording medium of domain wall displacement detection type 204 is read out by a magnetic head 206 and an optical head 205. The waveform distortion produced at this time is removed at the outset by a limiter circuit 201. After equalization by a PR equalizing circuit 202 to a PR (1, −1) channel, the waveform distortion is detected by a Viterbi detection circuit 203.    [Patent Document 1] JP Patent Kokai Publication No. JP-P2008-287763A, which corresponds to US Patent Application Publication No. US2010/0135142A1.    [Patent Document 2] JP Patent Kokai Publication No. JP-P2005-011385A, which corresponds to US Patent Application Publication No. US2004/0252589A1.    [Non-Patent Document 1] Ogawa, Homma et al., ‘Development of Technology of FID DVD Device Implementation (Recording Technique), Technical Report for Society of Information Media, ITE Technical Report, Vol. 28, No. 43, pp. 17 to 20, MMS2004-38, CE2004-39 (July 2004)"}
{"text": "It is well recognized that combating oil spills in coastal waters is of great communal interest. This has led to continuous efforts to improve oil spill mitigation strategies. The preferred spill mitigation strategy often uses oil booms to contain and concentrate floating oil, prior to oil skimmer recovery. Despite improvements over the decades, limitations of tow speed and operation under real world conditions remain serious.\nConventional Oil Boom Leakage\nA well-designed oil boom should be flexible to conform to wave motions, yet sufficiently rigid to retain as much oil as possible. Designs range from small, lightweight booms for manual harbor deployment, to large, robust booms for open sea use that need a crane and sizeable vessels to handle and deploy.\nImportant oil boom failure modes include overtopping where oil passes over the boom, and leakage where oil passes under the boom, surfacing as downstream oil patches. Boom overtopping mechanisms include overfilling, wave splashover, and boom-diving due to high towing speed. Boom leakage includes overfilling, frontal slick droplet injection (entrainment) from the confined pool, and boom-surfing due to opposing currents and/or winds blowing in the tow direction.\nOil pool thickness at the boom increases as more oil is collected and/or as the boom is towed faster. Overfilling often occurs with waves and current-induced boom diving and splashover. Leakage occurs when oil droplets are injected sufficiently deep to pass under the oil boom and is highly sensitive to towing speed. Once the towing velocity exceeds the boom's critical towing speed, the frontal wave ‘breaks’ on a large scale, increasing oil entrainment dramatically. Instability development and thus droplet injection increase strongly with towing speed.\nLeakage also occurs for acute angle flows at boom-segment junctions from vortices formation. These vortices can inject oil into the water that then underpasses the boom. Acute angle boom flows occur on the leading boom segments that steer oil towards the apex.\nBubble Oil Boom Background\nIn 1971, the U.S. Coast Guard tested a pneumatic boom for oil spill control, otherwise known as a bubble oil boom (BOB) (US Coast Guard. Heavy-duty oil containment systems: pneumatic barrier system. Report 714102/A/094, US Coast Guard Office of Research and Development, Contract DOT-CG-00-490-A to Oil Containment Division, Wilson Industries, 1971). A pneumatic boom generally is formed from a long, submerged air pipe with a series of holes along its length, typically at the upper generatrix. The curtain of bubbles rise in a sheet that drives an upwelling flow, which at the water surface is converted (by continuity) into an outwelling, which is the oil-blocking barrier. The Coast Guard study concluded that BOBs only were useful for low currents such as harbors, where current commercial BOBs are found. A number of other attempts at using BOBs have been made, such as in U.S. Pat. Nos. 3,491,023, 3,744,254, and more recently U.S. Patent Publication No. 2011/0303613. However, to date there are no large-scale commercial BOB systems available outside harbors due a number of drawbacks with previous designs."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a strip-line filter utilized to filter microwaves in a communication apparatus or a measuring apparatus operated in frequency bands ranging from an ultra high frequency (UHF) band to a super high frequency (SHF) band, and more particularly to a strip-line filter in which a strip line is shortened and is made plane at low cost. Also, the present invention relates generally to a dual mode resonator utilized for an oscillator or a strip-line filter, and more particularly to a dual mode resonator in which two types microwaves are independently resonated.\n2. Description of the Related Art\n2.1. First Previously Proposed Art:\nA strip-line resonating filter is manufactured by serially arranging a plurality of one-wavelength type of strip line ring resonators to reduce radiation loss of microwaves transmitting through a strip line of the resonating filter. However, there is a drawback in the strip-line resonating filter that the resonating filter cannot be downsized. Therefore, a dual mode strip-line filter in which microwaves in two orthogonal modes are resonated and filtered has been recently proposed. A conventional dual mode strip-line filter is described with reference to FIGS. 1 and 2.\nFIG. 1 is a plan view of a conventional dual mode strip-line filter. FIG. 2A is a sectional view taken generally along the line II--II of FIG. 1. FIG. 2B is another sectional view taken generally along the line II--II of FIG. 1 according to a modification.\nAs shown in FIG. 1, a conventional dual mode strip-line filter 11 comprises an input terminal 12 excited by microwaves, a one-wavelength strip line ring resonator 13 in which the microwaves are resonated, an input coupling capacitor 14 connecting the input terminal 12 and a coupling point A of the ring resonator 13 to couple the input terminal 12 excited by the microwaves to the ring resonator 13 in capactive coupling, an output terminal 15 which is excited by the microwaves resonated in the ring resonator 13, an output coupling capacitor 16 connecting the output terminal 15 and a coupling point B in the ring resonator 13 to couple the output terminal 15 to the ring resonator 13 in capactive coupling, a phase-shifting circuit 17 coupled to a coupling point C and a coupling point D of the ring resonator 13, a first coupling capacitor 18 for coupling a connecting terminal 20 of the phase-shifting circuit 17 to the coupling point C in capacitive coupling, and a second coupling capacitor 19 for coupling another connecting terminal 21 of the phase-shifting circuit 17 to the coupling point D in capacitive coupling.\nThe ring resonator 13 has a uniform line impedance and an electric length which is equivalent to a resonance wavelength .lambda..sub.o. In this specification, the electric length of a closed loop-shaped strip line such as the ring resonator 13 is expressed in an angular unit. For example, the electric length of the ring resonator 13 equivalent to the resonance wavelength .lambda..sub.o is called 360 degrees.\nThe input and output coupling capacitors 14, 16 and first and second coupling capacitors 18, 18 are respectively formed of a plate capacitor.\nThe coupling point B is spaced 90 degrees in the electric length (or a quarter-wave length of the microwaves) apart from the coupling point A. The coupling point C is spaced 180 degrees in the electric length (or a half-wave length of the microwaves) apart from the coupling point A. The coupling point D is spaced 180 degrees in the electric length apart from the coupling point B.\nThe phase-shifting circuit 17 is made of one or more passive or active elements such as a capacitor, an inductor, a strip line, an amplifier, a combination unit of those elements, or the like. A phase of the microwaves transferred to the phase-shifting circuit 17 shifts by a multiple of a half-wave length of the microwaves to produce phase-shift microwaves.\nAs shown in FIG. 2A, the ring resonator 13 comprises a strip conductive plate 22, a dielectric substrate 23 mounting the strip conductive plate 22, and a conductive substrate 24 mounting the dielectric substrate 23. That is, the ring resonator 13 is formed of a microstrip line. The wavelength of the microwaves depends on a relative dielectric constant .epsilon..sub.r of the dielectric substrate 23 so that the electric length of the ring resonator 13 depends on the relative dielectric constant .epsilon..sub.r.\nIn a modification, the ring resonator 13 is formed of a balanced strip line shown in FIG. 2B. As shown in FIG. 2B, the ring resonator 13 comprises a strip conductive plate 22m, a dielectric substrate 23m surrounding the strip conductive plate 22m, and a pair of conductive substrates 24m sandwiching the dielectric substrate 23m.\nIn the above configuration, when the input terminal 12 is excited by microwaves .lambda..sub.o, the electric field is induced around the input coupling capacitor 14 so that the intensity of the electric field at the coupling point A of the ring resonator 13 is increased to a maximum value. Therefore, the input terminal 12 is coupled to the ring resonator 13 in the capacitive coupling, and the microwaves are transferred from the input terminal 12 to the coupling point A of the ring resonator 13. Thereafter, the microwaves are circulated in the ring resonator 13 in clockwise and counterclockwise directions. In this case, the microwaves having the resonance wavelength .lambda..sub.c are selectively resonated according to a first resonance mode.\nThe intensity of the electric field induced by the microwaves resonated is minimized at the coupling point B spaced 90 degrees in the electric length apart from the coupling point A because the intensity of the electric field at the coupling point A is increased to the maximum value. Therefore, the microwaves are not directly transferred to the output terminal 15. Also, the intensity of the electric field is minimized at the coupling point D spaced 90 degrees in the electric length apart from the coupling point A so that the microwaves are not transferred from the coupling point D to the phase-shifting circuit 17. In contrast, because the coupling point C is spaced 180 degrees in the electric length apart from the coupling point A, the intensity of the electric field at the coupling point C is maximized, and the connecting terminal 20 is excited by the microwaves circulated in the ring resonator 13. Therefore, the microwaves are transferred from the coupling point C to the phase-shifting circuit 17 through the first coupling capacitor 18.\nIn the phase-shifting circuit 17, the phase of the microwaves shifts to produce phase-shift microwaves. For example, the phase of the microwaves shifts by a half-wave length thereof. Thereafter, the connecting terminal 21 is excited by the phase-shift microwaves, and the phase-shift microwaves are transferred to the coupling point D through the second coupling capacitor 19. Therefore, the intensity of the electric field at the coupling point D is increased to the maximum value. Thereafter, the phase-shift microwaves are circulated in the ring resonator 13 in the clockwise and counterclockwise directions so that the phase-shift microwaves are resonated according to a second resonance mode.\nThereafter, because the coupling point B is spaced 180 degrees in the electric length apart from the coupling point D, the intensity of the electric field is increased at the coupling point B. Therefore, electric field is induced around the output coupling capacitor 16, so that the output terminal 15 is coupled to the coupling point B in the capacitive coupling. Thereafter, the phase-shift microwaves are transferred from the coupling point B to the output terminal 15. In contrast, because the coupling points A, C are respectively spaced 90 degrees in the electric length apart from the coupling point D, the intensity of the electric field induced by the phase-shift microwaves is minimized at the coupling points A, C. Therefore, the phase-shift microwaves are transferred to neither the input terminal 12 nor the connecting terminal 20.\nAccordingly, the microwaves having the resonance wavelength .lambda..sub.o are selectively resonated in the ring resonator 13 and are transferred to the output terminal 15. Therefore, the conventional dual mode strip-line filter 11 functions as a resonator and filter.\nThe microwaves transferred from the input terminal 12 are initially resonated in the ring resonator 13 according to the first resonance mode, and the phase-shift microwaves are again resonated in the ring resonator 13 according to the second resonance mode. Also, the phase of the phase-shift microwaves shifts by 90 degrees as compared with the microwaves. Therefore, two orthogonal modes formed of the first resonance mode and the second resonance mode independently coexist in the ring resonator 13. Therefore, the conventional dual mode strip-line filter 11 functions as a two-stage filter.\n2.2. Problems of the First Previously Proposed Art to be Solved by the Invention:\nHowever, passband characteristics of the filter 11 is determined by the electric length of the ring resonator 13, so that a microwave having a fixed wavelength such as .lambda..sub.o is only resonated. Therefore, because the electric length of the ring resonator 13 is unadjustable, there is a drawback that the adjustment of the resonance wavelength is difficult.\nAlso, because it is required that the electric length of the strip line ring resonator 13 is equal to the one wavelength .lambda..sub.o of the resonance microwave and because the phase-shifting circuit 17 is formed of a concentrated constant element such as a coupling capacitor or a transmission line such as a strip line, there is another drawback that it is difficult to manufacture the filter 11 in a small-size and plane shape.\n2.3. Second Previously Proposed Art:\nFIG. 3 is a plan view of another conventional dual mode strip-line filter.\nAs shown in FIG. 3, another conventional dual mode strip-line filter 31 comprises two dual mode strip-line filters 11 arranged in series. An inter-stage coupling capacitor 32 is connected between the coupling point D of the filter 11 arranged at an upper stage and the coupling point A of the filter 11 arranged at a lower stage. The phase shifting circuit 17 of the filter 11 arranged at the upper stage is composed of a coupling capacitor 33, and the phase-shifting circuit 17 of the filter 11 arranged at the lower stage is composed of a coupling capacitor 34.\nIn the above configuration, when the input terminal 12 is excited by a signal (or a microwave) having a resonance wavelength .lambda..sub.o, the signal is resonated according to the first and second resonance modes in the same manner, and the signal is transferred to the coupling point A of the filter 11 arranged at the lower stage through the inter-stage coupling capacitor 32. Thereafter, the signal is again resonated according to the first and second resonance modes in the filter 11 arranged at the lower stage, and the signal is output from the coupling point D to the output terminal 15. In this case, the resonance wavelength .lambda..sub.o is determined according to an electric length of the ring resonator 13.\nTherefore, the conventional dual mode strip-line filter 31 functions as a four-stage filter in which the signal is resonated at four stages arranged in series.\n2.4. Problems of the Second Previously Proposed Art to be Solved by the Invention:\nHowever, it is required that the electric length of the strip line ring resonator 13 is equal to the one wavelength .lambda..sub.o of a resonance microwave, and it is required to increase the number of filters 11 for the purpose of improving attenuation characteristics of the resonance microwave. Therefore, there is a drawback that a small sized filter cannot be manufactured.\nAlso, the phase-shifting circuit 17 is formed of a concentrated constant element such as a coupling capacitor or a transmission line such as a strip line, there is another drawback that it is difficult to manufacture the filter 31 in a small-size and plane shape.\n2.5. Third Previously Proposed Art:\nA quarter-wavelength strip line resonator made of a balanced strip line or a micro-strip line has been broadly utilized in a high frequency band as an oscillator or a resonator utilized for a strip-line filter because the quarter-wavelength strip line resonator can be made in a small size. However, because ground processing in a high-frequency is performed for the quarter-wavelength strip line resonator, there are drawbacks that characteristics of a resonance frequency and a no-loaded Q factor (Q=.omega..sub.o /2.DELTA..omega., .omega..sub.o denotes a resonance angular frequency and .DELTA..omega. denotes a full width at half maximum) vary. To solve the drawbacks, a dual mode resonator in which two types microwaves having two different frequencies are resonated or a microwave is resonated in two stages by utilizing two independent resonance modes occurring in a ring-shaped resonator not grounded in high-frequency has been proposed for the purpose of downsizing a resonator. The dual mode resonator is, for example, written in a technical Report MW92-115 (1992-12) of Microwave Research in the Institute of Electronics, Information and Communication Engineers.\nA conventional dual mode resonator is described with reference to FIG. 4.\nFIG. 4 is an oblique view of a conventional dual mode resonator.\nAs shown in FIG. 4, a conventional dual mode resonator 41 comprises a rectangular-shaped strip line 42 for resonating two microwaves having two different frequencies f1 and f2, a lumped constant capacitor 43 connected to connecting points A, B of the rectangular-shaped strip line 42 for electromagnetically influencing the microwave having the frequency f1, a dielectric substrate 44 mounting the strip line 42, and a grounded conductive plate 45 mounting the dielectric substrate 44. Electric characteristics of the rectangular-shaped strip line 42 is the same as those of a ring-shaped strip line. The strip line 42 is made of a microstrip line. However, it is applicable that the strip line 42 be made of a balanced strip line.\nIn the above configuration, when a first input terminal (not shown) connected to the connecting point A is excited by a first signal (or a first microwave) having a frequency f1, an electric voltage at the connecting point A is increased to a maximum value. Therefore, the first signal is transferred from the first input terminal to the connecting point A of the strip line 42. Thereafter, the first signal is circulated in the strip line 42 in clockwise and counterclockwise directions in a first resonance mode. In this case, electric voltages at connecting points C and D spaced 90 degrees in the electric length (or a quarter-wave length of the first signal) apart from the connecting point A are respectively reduced to a minimum value, so that the first signal is not output from the connecting C or D to a terminal (not shown) connected to the connecting point C or D. Also, an electric voltage at the connecting point B spaced 180 degrees in the electric length (or a half-wave length of the first signal) apart from the connecting point A is increased to the maximum value, so that the first signal is output from the connecting point B to a first output terminal (not shown) connected to the connecting point B.\nIn contrast, when a second input terminal (not shown) connected to the connecting point C is excited by a second signal (or a second microwave) having a frequency f2, an electric voltage at the connecting point C is increased to a maximum value. Therefore, the second signal is transferred from the second input terminal to the connecting point C of the strip line 42. Thereafter, the second signal is circulated in the strip line 42 in clockwise and counterclockwise directions in a second resonance mode. In this case, electric voltages at the connecting points A and B spaced 90 degrees in the electric length apart from the connecting point C are respectively reduced to a minimum value, so that the second signal is not output from the connecting point A or B to the first input or output terminal connected to the connecting point A or B. Also, an electric voltage at the connecting point D spaced 180 degrees in the electric length apart from the connecting point C is increased to the maximum value, so that the second signal is output from the connecting point B to a second output terminal (not shown) connected to the connecting point D.\nBecause any lumped constant capacitor connected to the connecting points C and D is not provided, the frequency f1 differs from the frequency f2. However, in cases where a capacitor having the same capacity as that of the capacitor 43 is provided to be connected between the connecting points C and D, the frequency f2 is equal to the frequency f1. Also, in cases where the capacitor 43 is removed, the frequency f1 is equal to the frequency f2. Therefore, the frequencies f1 and f2 resonated in the first and second resonance modes independent each other are the same. In other words, the conventional dual mode resonator 41 functions as a two-stage resonator in which two microwaves having the same frequency are resonated in two stages arranged in parallel.\nAccordingly, the resonator 41 comprising the strip line 42 and the capacitor 43 functions as a dual mode resonator in which two microwaves are resonated in two resonance modes independent each other. Because the resonator 41 is not grounded in high-frequency as a special feature of a dual mode resonator and because radiation loss of the microwave is lessened because of a closed-shape strip line as another special feature of the dual mode resonator, the resonator 41 can be manufactured in a small size without losing the special features of a one-wavelength ring-shaped dual mode resonator.\n2.6. Problems of the Third Previously Proposed Art to be Solved by the Invention:\nHowever, it is required to accurately set a lumped capacity of the capacitor 43 for the purpose of obtaining a resonance frequency of a microwave at a good reproductivity. In actual manufacturing of the dual mode resonator 41, it is difficult to accurately set a lumped capacity of the capacitor 43. In cases where a frequency adjusting element is additionally provided for the dual mode resonator 41 to accurately set a lumped capacity of the capacitor 43, the number of constitutional parts of the dual mode resonator 41 is increased. Therefore, there are drawbacks that resonating functions of the resonator 41 are degraded and a manufacturing cost of the resonator 41 is increased."}
{"text": "The present invention relates generally to enclosures for switchgear, panel boards, circuit breakers, and more particularly to a non-walk-in enclosure for low voltage switchgear and switchboard assemblies.\nSwitchgear and switchboard are general terms which cover metal enclosures, housing switching and interrupting devices such as fuses, circuit breakers, relays along with associated control, instrumentation and metering devices. The enclosures also typically include devices such as bus bars, inner connections and supporting structures used for the distribution of electrical power. Such metal enclosures can be maintained in a building such as a factory or commercial establishment, or they can be maintained outside of such facilities and exposed to environmental weather conditions. Typically, hinge doors or covers are provided on the front of the switchgear or switchboard sections for access to the devices contained therein.\nLow voltage switchgear and switchboards operate at voltages up to 600 volts and continuous currents that can exceed 5000 amps. It is also possible that in short circuit conditions, intermittent currents exceeding 100,000 amps are possible.\nWhen metal enclosures are used out of doors, switchgear and switchboard assemblies are typically housed in NEMA type 3R enclosures. Previous non-walk-in enclosures had the exterior doors closed over the circuit breaker doors and sealed against the indoor switchgear frame when the frame was wider than the circuit breaker doors. Other enclosures had the door sealed against a spacer that was used in the same plane as the front of the switchgear frame when the circuit breaker doors and the frame are nominally the same width. The spacer was needed to provide a sealing surface for the exterior door. Such prior art spacers are typically approximately two inches wide and are placed between each section of switchgear and switchboards. Such spacers add cost and proliferates parts to accommodate the added width of the enclosure. Conventional enclosure designs also lack space between the circuit breaker door and the exterior door of the enclosure thereby preventing a draw-out type circuit breaker from being placed in a disconnected position with the exterior door of the enclosure closed. In a situation where the enclosure is housed inside of a building, the additional width requirement takes up additional space.\nThus, there is a need for a non-walk-in enclosure for low voltage switchgear and switchboard assemblies that does not require the additional width of spacers between sections. There is a further need for an enclosure that will accommodate circuit breaker compartment doors that are nominally the same width as the sections in which the circuit breakers are housed. There is an additional need for an enclosure to allow the exterior doors to be closed when an enclosed circuit breaker is in a disconnected position."}
{"text": "1. Field of the Invention\nThis invention relates to a flatness measuring apparatus for measuring flatness of a surface of a plate-shaped object such as a magnetic disk to be used as storage medium for computers.\n2. Related Art Statement\nFor example, magnetic disks to be used as storage medium for computers typically comprise a substrate made of aluminum alloy or glass and a magnetic film formed on the surface of the substrate.\nFor manufacturing high quality magnetic disks, the disks are required to show a certain level of surface flatness, so their substrates are normally checked for flatness of the surface in an after-fabrication process.\nElectric micrometers are known as means for checking whether the object of examination such as a plate-shaped product shows a predetermined level of surface flatness. The electric micrometer is in fact a displacement sensor comprising a probe and a transducer for transducing any mechanical displacement of the probe into an electric signal such that, as the probe is driven to slidingly move on the surface of a plate-shaped object, any displacement of the probe on the surface is transferred to the core of the actuator transformer of the sensor and detected as electric signal representing the surface condition of the object.\nHowever, with a method of measuring the surface flatness of an object by holding a contact probe in contact with the surface of the object, the contact probe may damage and leave one or more than one flaws on the surface. Such a flaw should be avoided by any means particularly when it comes to a magnetic disk that requires a mirror finish to say nothing of a highly smoothed surface. Additionally, for the purpose of measuring efficiency, the operation of observing the surfaces of products supplied on a mass production basis is preferably automated."}
{"text": "1. Field of the Invention\nThe present invention relates generally to solder masks for use on carrier substrates and, more particularly, to solder masks that are configured to prevent excess adhesive material from contacting or covering terminals of carrier substrates. In addition, the present invention relates to carrier substrates including such solder masks and to semiconductor device assemblies and packages including the carrier substrates. The present invention also relates to methods for designing solder masks to be used on carrier substrates, to methods for forming solder masks on carrier substrates, and to methods for connecting semiconductor devices to carrier substrates.\n2. Background of Related Art\nCircuit boards and other carrier substrates have long been employed to establish electrical communication between bond pads of semiconductor dice and other electronic components, including those carried by the circuit board and components that are external to the circuit board.\nConventionally, relatively small circuit boards have been used to package one or more semiconductor dice, which protects the dice and provides a semiconductor device package that includes a standardized connection pattern, regardless of the arrangement of bond pads of each semiconductor die of such a package. These circuit boards are often referred to in the art as “interposers.” In the state of the art, the dimensions of interposers are approaching the corresponding dimensions of the semiconductor dice to be secured thereto, with the result being a continuous decrease in the sizes of semiconductor device packages, many of which are so-called “chip-scale packages” exhibiting dimensions that are the same as or only slightly larger than those of the die thereof.\nNonetheless, when conventional wire bonding processes are used to electrically connect bond pads of a semiconductor die to corresponding contacts, or terminals, of an interposer to which the semiconductor die is back bonded, the terminals of the interposer must be positioned outside of the outer periphery of the semiconductor die. This requirement prevents a reduction in the size of these interposers to dimensions that are substantially the same as those of the semiconductor die to be secured thereto and carried thereby.\nMoreover, in order to securely attach a semiconductor die to the circuit board, it is typically desirable that a sufficient amount of adhesive material be used to form a confluent adhesive layer adjacent to the back side of the semiconductor die and a fillet around the entire outer periphery of the semiconductor die. Some excess adhesive material is typically required to adhere the semiconductor die and the circuit board to one another in this manner. When excess adhesive material is used, however, some of the adhesive material may flow outward from the semiconductor die, beyond the desired location for a fillet, and contaminate the connection surfaces of adjacent terminals, which are located at the same elevation as or just above a surface of the circuit board. If the connection surfaces of the terminals become contaminated in this manner, intermediate conductive elements, such as bond wires, may not adhere thereto or form an effective electrical connection, which may ultimately reduce the reliability of the packaged semiconductor device, or cause electrical shorting and failure of the packaged semiconductor device.\nLarger circuit boards, such as so-called “motherboards” and the substrates of so-called “plug-and-play” devices or cards, may be configured to carry a plurality of semiconductor devices. As with interposers, the trend with such larger circuit boards is also toward ever-decreasing sizes and ever-increasing densities. Consequently, contamination of the connection surfaces of terminals by adhesive material is increasingly encountered when semiconductor devices are assembled with larger circuit boards.\nSolder masks are often used on circuit boards to facilitate the formation of solder balls or bumps on selected terminals of the circuit boards while protecting other features thereof and preventing solder from wetting and adhering to other features of the circuit board. The solder masks that are used on circuit boards are typically fabricated from photoresists, using known photolithography processes. These solder masks may be left in place while semiconductor devices and other components are secured and electrically connected to the circuit board. Thus, a semiconductor device may be secured to a solder mask which is, in turn, secured to the circuit board to which the semiconductor device is to be electrically connected.\nWhen conventional solder masks are used on conventional circuit boards, a portion of the adhesive material that secures a semiconductor device to the solder mask may still flow laterally outward, beyond the outer periphery of the semiconductor device and onto connection surfaces of the adjacent terminals of the circuit board.\nSome solder masks have been configured with recesses to receive some of the excess adhesive material. Nonetheless, these recesses typically do not have great enough volumes to contain all of the excess adhesive material that is forced from beneath a semiconductor device as the semiconductor device is secured to the solder mask and, thus, to the circuit board. Consequently, the terminals of circuit boards are typically spaced a sufficient distance away from the location at which the outer periphery of a semiconductor device to be secured to the circuit board is to be positioned that the connection surfaces thereof will not be contaminated by adhesive material. Such positioning of the terminals undesirably requires that the circuit board include additional real estate, which prevents a reduction in the dimensions of the circuit board.\nThe inventors are not aware of any carrier substrates or solder masks therefor that are configured to prevent adhesive material from flowing onto the connection surfaces of the carrier substrate terminals while reducing the dimensions of the carrier substrate."}
{"text": "This invention relates to gas detecting apparatus, and more particularly, to natural gas detecting apparatus which is able to discriminate between pipeline natural gas, non-pipeline sources of methane, propane, and gasoline vapors.\nUtilities which distribute natural gas require reliable gas-leak detectors for use in maintenance of gas supply lines. Existing natural gas detectors are either costly, sensitive and non-selective or low cost, insensitive and non-selective. Non-selective gas detectors respond to any combustible gas. Selective gas detectors are specific to hydrocarbon gases. The two presently most used gas detectors are based on hydrogen flame ionization and on hot wire catalysis. These gas detectors cannot distinguish among different types of hydrocarbons. However, it is necessary to distinguish among different types of hydrocarbons in order to distinguish a pipeline gas from gasoline vapors or sewer or swamp gas and so reduce leak surveyor time wasted on false alarms. Ethane content, if measurable, provides a good means to discriminate between pipeline gas and interfering gasoline vapors and sewer or swamp gases because the later contain practically no ethane, while pipeline gas does, in varying degrees. Gasoline vapors and propane (LP gas) can also generate a false alarm with conventional instruments. However, their infrared absorption is shifted relative to that of methane, as will be described later, as is the basis for this invention to eliminate such false alarms.\nIn U.S. Pat. No. 4,507,558, there is disclosed a selective detector for natural gas which discriminates between low concentrations of natural gas and other methane sources by measuring the characteristics of the methane/ethane ratio of natural gas as well as by using a combustible gas sensor. The operation of this detector is based on infrared light absorption of methane and ethane in combination with another non-specific combustible gas detector whereby the detector has the ability to detect nonspecifically, the presence of a combustible gas, and to define the nature of the combustible gas. Thus, this natural gas detector utilizes two types of detection including nondispersive infrared detectors and a non-specific combustible detectors such as hot-wire catalytic combustible detector. The detector determines concentration of both methane or ethane irrespective of the concentration of the other gas by using absorption cells placed in front of the detectors. The detector includes a light emitting diode which issues light centered around 3.32 microns and a reference light source which emits light at a wavelength outside of this band. Although this arrangement permits distinguishing among different types of hydrocarbons, the requirement for a hotwire catalytic combustible detector adds cost and complexity to the device and increases power consumption.\nIt would be desirable to have a natural gas detector which can distinguish among different types of hydrocarbons, and which provides information to the user on the amount and type of combustible gases in the environment."}
{"text": "A system on a chip (SOC) generally includes at least one processing core, which generally is operatively coupled to a level 2 (L2) memory cache. A cache in a SOC typically needs to fetch, from a memory, current instructions and current data (as and when required by the processing core, in case such current instructions and current data are, for example, not already cached in the cache and/or is dirty), as well as pre-fetch instructions and pre-fetch data corresponding to instructions and data that are likely to be needed, by the processing core, in a forthcoming operation. In conventional SOC architectures each of the current instructions, current data, pre-fetch instructions and pre-fetch data are communicated between the processor and the cache, such as an L2 cache, via dedicated ports.\nThe description in this section is related art, and does not necessarily include information disclosed under 37 C.F.R. 1.97 and 37 C.F.R. 1.98. Unless specifically denoted as prior art, it is not admitted that any description of related art is prior art."}
{"text": "The invention relates to a method of manufacturing a large-sized, thin-walled object of fiber reinforced synthetic resin using inner and outer mold parts which bound a cavity for forming the aforementioned thin-walled object, in which a reduced pressure is created in the cavity by means of at least one vent in one of the two mold parts, which vent communicates with a vacuum source, and then flowing a hardening synthetic resin into the cavity via one or more openings in at least one of the two mold parts.\nSuch a method is known from applicant's Dutch patent specification 178,761 and the corresponding U.S. patent specification 4,359,437.\nThe disadvantage of this method is that in addition to a form-retaining inner mold part, a flexible, form-retaining outer mold part is also required. This means that, for forming large objects of different dimensions, different mold parts are required.\nThe invention overcomes the aforementioned disadvantage."}
{"text": "The invention relates to a coolable trunk piston for internal combustion engines.\nSuch pistons are known from U.S. Pat. No. 4,377,967. The wall part (shaped sheet metal) which is used in those pistons to cover the cavity which forms the cooling oil space is flanged onto the outer periphery of the annular wall which carries the piston rings. This type of fixing is not secure, because flanging of the relatively thin-gauge wall part can become detached in the long term under the conditions in which engines operate.\nThe use of a one-piece thick gauge and thus, by reason of its greater rigidity, a wall part which guarantees a more reliable attachment is in turn undesirable for different reasons. The outer annular part of the piston head which accommodates the piston rings ought during operation of the engine to be able to change its shape at its bottom end without being hindered by the incorporated wall part. Such deformations take place due to the combustion processes which take place in the manner of explosions and, if they are inhibited, they can cause cracks in the piston head material. In this respect, one-piece sheet metal covering such as are known from FR-A 23 23 022 or coverings which consist of individual parts rigidly connected to each other, such as are described in the case of pistons in U.S. Pat. No. 2,698,210 are equally unsuitable for the necessary covering of the annular cooling oil cavity.\nIt is on this premise that the invention is based on the problem in the case of the piston known from U.S. Pat. No. 4,377,967, of providing a covering which ensures on the one hand movements of the bottom end of the outer annular wall of the piston head while on the other is sufficiently securely attached to the piston head."}
{"text": "Gangliosides are sialic acid containing glycosphingolipids composed of a complex carbohydrate moiety linked to a hydrophobic ceramide portion. Embedded within the outer leaflet of the cell membrane, the carbohydrate chain is exposed to the extracellular matrix. Qualitative and quantitative changes in ganglioside composition during cell differentiation and proliferation have been observed and seem to reflect the state of malignant transformation of cancers of neuroectodermal origin (Hakomori 1985). Malignant melanoma cells express a variety of complex gangliosides in addition to GM3, the major ganglioside in normal melanoctyes (Carubia et al. 1984). Altered ganglioside metabolism in melanoma causes additional expression of GD3, GD2, GM2,9-O-Acetyl-GD3 and GT3 (Hamilton et al. 1993; Tsuchida et al., 1987). Treatment of patients with anti-GD3 monoclonal antibodies resulted in inflammation at the tumor site and partial regression of metastasis was seen occasionally, suggesting, that gangliosides are suitable targets for immune attack (Houghton et al., 1985). The generation of human MAb's reactive with GD3 from melanoma patients (Yamaguchi et al., 1987) support the idea, that gangliosides are potential immunogens as well.\nIn studies aimed at inducing a humoral response against gangliosides in melanoma patients by active immunization, GM2/BCG vaccines seemed to be most effective (Livingston et al., 1987; Livingston et al., 1989). In a randomized study with 122 melanoma patients, who were disease-free after surgery, that it was showed that, out of 64 patients treated with BCG alone and 58 patients with GM2/BCG, the majority of patients (86%) receiving the GM2 vaccine produced antibodies. Patients that produced anti-GM2 antibodies had a significantly longer disease free and overall survival than antibody negative patients. Comparing the two arms of the trial, patients receiving the GM2/BCG vaccine had a 17% improvement in disease-free interval and 9% improvement in survival when compared to the BCG control group, though neither result was statistically significant (Livingston et al., 1993a). Unfortunately, the immune response was only of short duration, mostly IgM and of moderate titer. This suggested that GM2 was recognized as a T-cell independent antigen as a consequence of carbohydrate antigens (Livingston et al., 1989) and also because gangliosides are auto antigens expressed on some normal tissue (Hamilton et al., 1993). Similar approaches with GD2 and 9-O-Acetyl-GD3 vaccines in patients resulted in occasionally low titers and no antibody response against GD3 could be detected (Livingston, 1991).\nNew potent adjuvants were able to enhance the immune responses against gangliosides in some cases, but especially for auto antigens such as for tumor associated gangliosides a different approach had to be utilized. Based on Landsteiner's classical experiments (Landsteiner and Chase, 1942) with hapten-carrier conjugates, covalent attachment of poorly immunogenic antigens to immunogenic carrier proteins has been successfully used to enhance immune response. For example responsiveness to carbohydrates, other than gangliosides, could be accomplished with conjugation to appropriate carrier proteins. Coupling of bacterial capsular polysaccharides to immunogenic proteins showed a significant increase in immune response and protection (Eskola et al., 1990). Recently, vaccination of ovarian cancer patients with synthetic Thompson Friedenreich tumor antigen conjugated to keyhole limpet hemocyanin elicited humoral IgM and IgG response (MacLean et al., 1992). The important finding common in these studies was the isotype switch from a IgM response of short duration to a long lasting, high affinity IgG response indicating that activation of T-cell dependent pathways against carbohydrates is likely to occur. This approach is now applied to the melanoma tumor antigen GD3 to develop a method to synthesize ganglioside-protein conjugate vaccines and examine the immunogenicity of different GD3-protein conjugates in mice."}
{"text": "1. The Field of the Invention\nThe present invention relates to wireless communication networks and related systems and devices. More specifically, embodiments of the present invention relate to systems and methods for using a programmable agent within a wireless device within a distributed communications network to monitor wireless networks, devices, and/or the operation thereof.\n2. The Relevant Technology\nAs a result of advances in technology and enormous increases in the number of wireless device users, the size and complexity of wireless communications networks has greatly increased. One consequence of such increases in size and complexity is that the relative increase in operational and performance problems associated with communications networks has also increased. Reliability issues, such as dropped calls, lack of coverage, poor audio quality, and application failure often lead to user frustration and to increased costs. As new services are introduced that use even more complex technology, exercise different usage modalities, and place additional demands on networks, network performance continues to be a prime concern. In fact, quality of service often has a direct impact on a service provider's profitability. Therefore, improving quality of service is a top priority for service providers.\nNetwork monitoring solutions are well known in the art and widely employed by service providers. Unfortunately, the solutions that are currently available can only monitor and diagnose subsets of the overall telecommunications system in limited ways rather than providing a holistic view of network and device performance that may be needed to efficiently identify and resolve quality issues.\nTypical approaches to network monitoring include “self-monitoring.” In this case, a network element reports on its own operational status, performance and any errors that occur during its operation. The metrics reported by a single element can sometimes be indicative of a broader, system-wide problem, but rather than providing answers, problem resolution entails guesswork and extended troubleshooting, which wastes valuable resources.\nAnother common approach includes placing probes at various points in the network to determine if network elements are functioning according to specification. Sometimes referred to as “sniffers”, “log monitors” or “event monitors,” these monitoring systems are effective at identifying performance issues with a particular network element, but they fail to capture problems that stem from the interfaces among network elements, i.e., these solutions do not address the case where single elements are performing according to specifications, but problems occur when those elements interact with one another. This far more complex and subtle set of problems has costly consequences to network operators when services cannot be delivered to end customers.\nAnother monitoring approach involves pre-programmed service monitors, where specific elements perform service transactions to emulate “real-world” transaction activity; end to end performance is then monitored and the results reported. While these solutions may catch systematic failures, they cannot detect intermittent or dispersed problems, subtle impairments, or device or end user specific issues. Further, they can only test anticipated usage scenarios and fail to adapt to new usages and interactions between services.\nSignificantly, the aforementioned solutions lack the ability to monitor network conditions and intelligently and dynamically define and generate data reporting models in response to changing network conditions. Even with the employment of probes in a communications network, it is often the case that the probe provides only an indication of a problem, and actually troubleshooting the problem requires personnel to be dispatched to a physical location on the network, adding significant time and cost to problem identification and resolution. Moreover, these troubleshooting techniques do not provide a multi-faceted view of the different network layers, such as the physical layer, the applications layer, etc., and they do not correlate performance issues across these layers. As a result, numerous quality issues impacting end users go undetected or are misunderstood. Consequently, they may become crises because the performance data provided by currently available network monitoring solutions lacks the detail and timeliness needed to quickly identify, prioritize and resolve network issues.\nFurthermore, currently available network monitoring solutions are unable to adequately monitor and report a particular end user's experience, meaning that the service providers must rely upon the end user to report performance problems to a customer service representative. As such, it is frequently the case that by the time the user has reported a problem, the problem cannot be duplicated due to the user's difficulty recounting the specific details of what the user experienced, the timing of the occurrence, along with a lack of underlying data to validate the information provided by the user. Additionally, because the service provider is unable to view network performance holistically, it is difficult for the service provider to understand the true scope of the reported problem and appropriately prioritize the problem for resolution. Therefore, the problem may not be resolved to the reporting end user's satisfaction, creating a disincentive for the end user to diligently report problems in the future. Furthermore, in such situations, the opportunity to prevent the issue from affecting other users of the system and to improve the overall quality of the network is missed.\nThus, situations frequently arise where the end user is alienated from the company providing the communications service, without the service provider even being aware of the cause or source of the end user's dissatisfaction. Moreover, because so many end user problems go unreported, it is almost impossible for service providers to obtain reliable information about the scope of a particular issue. Therefore resources are frequently spent solving issues that are well described, but only affect a small number of end users, versus problems that are poorly described but effect a much larger number of end users. As a result, a service provider with an incorrect understanding of the scope of an issue may allocate an inappropriate amount of resources toward resolution, or it may fail to recognize the value of prioritizing resolution.\nOther monitoring solutions known in the art are directed to gathering and analyzing information about the performance of wireless devices. Typically, the monitoring software will be installed on a wireless device at the time of manufacture or by downloading the software onto the device. The software then runs continuously in the background, monitoring device and application performance. When a particular event or error associated with the device occurs, the software collects the data associated with the error or event and may upload it either in real time or at a later time to a data repository for analysis. One problem with such systems and methods, however, is that the systems require a relatively large amount of data to be uploaded to the data repository for analysis, which can drain the resources of the wireless device. Specifically, the frequent transfer of data or the transfer of large amounts of data may decrease the battery life and impact processing performance of the wireless device, which may create further user frustration and dissatisfaction with the communication service."}
{"text": "The present invention relates to oxygen control systems and more particularly to an oxygen control system for controlling the level of oxygen in fresh food storage containers.\nThe advent of containerized refrigerated cargo transport has allowed fresh agricultural products to be shipped to ever more distant markets and has enabled the transfer of containers between different modes of transport (e.g., truck to ship to truck) as required without rehandling of the cargo. Insect infestation of raw agricultural products is a universal problem in produce storage, processing and shipping that must be effectively addressed to minimize health concerns and to maintain an appetizing appearance, shelf life and marketability. Cargos that cross national boundaries are also generally subject to stringent regulations for all the same reasons as well as to prevent the importation of harmful new agricultural pests. Compared to losses in the field, the losses become increasingly significant during transportation as further value is added to the product in harvesting, packaging and shipping. Disinfestation can be achieved by heating or freezing but not without some sacrifice in quality, and the use of insecticides and fungicides is becoming increasingly unacceptable.\nThe use of containerized cargo transport in which containers carry a dedicated refrigeration unit (\"reefers\") has presented the opportunity to employ controlled-atmospheres for preservation of perishable products. Controlled-atmosphere (CA) shipment of food products and other perishables for preservation has received significant attention in recent years. For example, Transfresh Corporation has developed technology to effectively \"seal\" reefers and to sufficiently control the oxygen and other gas composition within the container to extend the fresh shelf-life of many agricultural and perishable products. Recently, systems have been developed for active continuous removal of oxygen leaking into \"sealed\" reefers to very low levels (&lt;1000 ppm) to extend CA technology to the shipment of fresh meat and fish products.\nThe principle of hermetic storage has been known from pre-industrial times and has been widely implemented for grain, bean and oil-seed storage in underground pits; the first large scale tests were run in Paris from 1819 to 1830. This principle has been steadily improved upon and extended ever since. Existing oxygen removal systems for shipment of fresh meat and fish are on-board systems that actively maintain oxygen at the lowest level possible, compensating for any air leakage throughout the shipping period (2-3 weeks). On-board systems need to be small for the least product displacement and for conservation of power; this is achieved by pre-flushing the container, tight sealing against leaks and pressure control to reduce the initial power demand and the subsequent load on the system. This very low concentration of oxygen would generally be lethal to most insects but it is not generally compatible with produce that respires. Iceberg lettuce, for example, requires &gt;1200 ppm of oxygen to prevent spoilage, but, &lt;1800 ppm at 36.degree. F. over a three-day exposure to kill green peach aphids; this is followed by 10 days at 2% O.sub.2 in transit, the normal CA shipping environment. The requirements for a disinfestation system then would be:\n1) Controlled oxygen levels rather than minimum oxygen maintenance. PA1 2) Continuous oxygen concentration measurement of sufficient resolution to set the required oxygen levels for a variety of products.\nThe specific technology solutions that have been used to address the above problems will now be discussed.\na. Modified-Atmosphere Disinfestation.\nHermetic storage has long been used for disinfestation of grain, beans and oil-seeds stored in underground pits. This process involves sealing the product in a pit, bin or silo and allowing the respiration of the stored product to reduce the oxygen to a level lethal to insects. Primitive versions of hermetic storage are in active use in many third world countries with mild winters, and more modern concrete storage containers are in continuous use for corn storage in Cyprus and Kenya to provide famine protection. This concept of modifying atmospheres was extended in the 1860s by adding nitrogen or \"burned air\" to grain storage silos. However serious interest in using the technique in a practical, routine manner was not pursued until the 1950s and 1960s. Until this time conventional fumigants and grain protectants had been successful in controlling stored-product pests. At that time, people began to realize that the chemicals, if used improperly, left objectional residues and were hazardous to apply. There was also a potential for the development of insect resistance to these chemicals. Research and development in controlled atmosphere disinfestation has been on-going since the 1960s with considerable success. It has been found that low temperatures, low oxygen and high carbon dioxide separately or in various combinations can be effective for disinfestation. For fresh produce there is the additional requirement to maintain sufficient oxygen to prevent acceleration of ripening or spoilage. The conditions of atmosphere and exposure time to achieve disinfestation without excessive ripening tends to be complex and unique to the combination of agricultural product and indigenous pest(s).\nA method and apparatus for storing foodstuffs using an electrolytic cell is described in U.S. Pat. No. 4,212,891 issued to Fujita and Kudo. This apparatus included a liquid alkaline electrolyte electrolytic stack and system, but did not include integral reactor feedback sensing. Moreover, with an alkaline system, the ambient CO.sub.2 reacts with the electrolyte and frequently there is liquid electrolyte leakage into the seal and manifold area, leading to competing reactions and shunt currents. Also, at near-ambient temperatures, water and electrolyte concentration management is difficult with a liquid electrolyte system. This makes feedback sensing and control difficult, especially if the same stack would be used for oxygen depletion, monitoring and control.\nb. Controlled-Atmosphere (CA) Technology.\nThe Tectrol.RTM. CA systems developed by Transfresh and in wide use (more than 22,000 CA-equipped reefers worldwide) are implemented in three stages: 1) sealing of specially equipped containers after loading; 2) inert-gas flushing down to 1-2% oxygen; and 3) controlled leakage of air and removal of carbon dioxide to compensate for respiration, based on regular oxygen and carbon dioxide concentration measurements during shipment. The addition of oxygen is accomplished with a unique air-exchange system; a slide valve inside the Tectrol controller, which is mounted in an access panel at the front of the reefer, allows precise amounts of pressurized air generated by the reefer fan to be pushed out through an air-exchange outlet port and outside air to be drawn into the low pressure side of the fan through an air-exchange inlet port. The removal of carbon dioxide is accomplished by periodically circulating the reefer atmosphere through a lime bed based on carbon dioxide sensor measurements. The atmosphere thus achieved lowers respiration extending the fresh life of the produce over a longer period of time but it is not typically adequate to effect disinfestation.\nIt is therefore a principal object of the present invention to provide a proton-exchange-membrane based electrochemical oxygen control system that can deplete, monitor and control oxygen in a closed container to levels sufficient to achieve disinfestation.\nIt is another object to the present invention to provide a proton-exchange-membrane based electrochemical oxygen control system that can provide feedback sensing to control oxygen levels using the same electrochemical cells stack for oxygen depletion, monitoring and control."}
{"text": "The present invention relates to a resin composition comprising a boron-containing polymer and in particular to a resin composition comprising a boron-containing polymer useful as an antifouling coating.\nAquatic creatures such as barnacles, ascidians, pylopagurus serpulophilus, blue mussels, bugula neritina, green layer and sea lettuce adhere to various substances below the water surface, to cause a wide variety of damage. For example, when these aquatic creatures adhere to the body of a ship, the velocity of the ship is lowered and the fuel cost is increased. Further, it is known that when the aquatic creatures adhere to fishing nets etc., the net maybe clogged to cause death of fish.\nFor preventing adhesion and growth of aquatic creatures on substances below the water surface, an antifouling coating is used, and there has been used an anti fouling coating comprising seawater-insoluble vinyl resin and alkyd resin as binder resin and a trialkyl tin polymeric compound as a vehicle. However, there is anxiety about the adverse effect of the eluted trialkyl tin on the ecosystem, so there is demand for development of heavy metal-free antifouling coatings and antifouling agents.\nJapanese Kokai Publication Hei-7-133207 discloses triphenyl borane-pyridine derivative complex compounds. Further, triphenyl borane-fatty amine complex compounds (Japanese Kokai Publication Hei-8-295608, Japanese Kokai Publication Hei-8-295609 etc.), triphenyl borane-hydroxyl group-containing amine complex compounds (Japanese Kokai Publication Hei-9-78007), triphenyl borane-higher fatty secondary or tertiary amine complex compounds (Japanese Kokai Publication Hei-10-182322) etc. are also proposed.\nFor the purpose of expressing an antifouling effect for a prolonged period of time with good coating workability, Japanese Kokai Publication Hei-11-322763 discloses triphenyl borane-alkoxyl group-containing amine complex compounds, and Japanese Kokai Publication Hei-11-323206 discloses triphenyl borane-(iso)nicotinate complex compounds.\nFurther, Japanese Kokai Publication Hei-11-199801 discloses, as a resin composition for heavy metal-free antifouling coatings, an antifouling coating composition comprising N atom-coordinated triaryl borane-containing polymer as a vehicle.\nThe object of the present invention is to provide a resin composition comprising a boron-containing polymer useful as an antifouling coating which does not contain heavy metal from an environmental viewpoint.\nThus, the present invention provides a resin composition\nwhich comprises a boron-containing polymer having, in a side chain or at a terminal thereof, at least one group represented by the following formula (1): \nin the formula, R1 is a part of an acid group; R2 groups are the same or different and each represents a substituted or unsubstituted alkyl or aryl; X represents an oxygen or sulfur atom; Y represents a primary to tertiary amine or a heterocyclic compound containing a substituted or unsubstituted nitrogen atom; and n is 1 or 2.\nPreferably, the above R1 is represented by the following formula (2). \nFurther, the present invention provides a production method of a boron-containing polymer\nwhich comprises a step of reacting a polymer having at least one acid group in a side chain or at a terminal thereof, with a nitrogen atom-coordinated boron compound represented by the following formula (3): \nin the formula, R2 groups are the same or different and each represents a substituted or unsubstituted alkyl or aryl; and Y represents a primary to tertiary amine or a heterocyclic compound containing a substituted or unsubstituted nitrogen atom.\nPreferably, the above acid group is a carboxylic acid group, a sulfonic acid group or a phosphoric acid group.\nFurthermore, the present invention provides a production method of a boron-containing polymer\nwhich comprises a step (I) of reacting a polymerizable organic acid with a nitrogen atom-coordinated boron compound represented by the following formula (3): \nin the formula, R2 groups are the same or different and each represents a substituted or unsubstituted alkyl or aryl; and Y represents a primary to tertiary amine or a heterocyclic compound containing a substituted or unsubstituted nitrogen atom,\nand a step (II) of polymerizing a monomer component comprising a boron-containing polymerizable monomer obtained in said step (I).\nPreferably, said polymerizable organic acid is a polymerizable organic carboxylic acid, a polymerizable organic sulfonic acid or a polymerizable organic phosphoric acid.\nMoreover, the present invention provides a resin composition\nwhich comprises a boron-containing polymer obtainable by the above production method of the boron-containing polymer.\nMoreover, the present invention provides an antifouling coating\nwhich comprises the above resin composition.\nThe resin composition of the present invention comprises a boron-containing polymer having at least one group represented by the above formula (1) in a side chain or at a terminal thereof.\nIn the above formula (1), R1 represents a part of an acid group. The above R1 is preferably a part of carboxylic acid group, sulfonic acid group or phosphoric acid group represented by the above formula (2). The R1 is more preferably a part of carboxylic acid group.\nThe above R1 constitutes an acid group in the form of xe2x80x94R1 (Xxe2x80x94H)n.\nThe above X represents an oxygen or sulfur atom. It is preferably an oxygen atom.\nIn the above formula (1), n is 1 or 2. When a part of acid group as the above R1 is a part of phosphoric acid group, n can be 2.\nIn the above formula (1), R2 groups are the same or different and each represents a substituted or unsubstituted alkyl or aryl. The above alkyl group may be linear or branched.\nThe number of carbon atoms in the alkyl or aryl group is preferably 1 to 30. If the number of carbon atoms is 30 or more, the antifouling property will be deteriorated when the composition is incorporated into an antifouling coating so that it will be necessary to be incorporated in a large amount, in some cases. Further, the resulting resin composition may be hardly dissolved in an organic solvent and is inferior in handling properties. The number of carbon atoms therein is more preferably 1 to 20.\nThe substituent group on the above alkyl or aryl includes an alkyl group, carboxyl group, sulfonic acid group, phosphoric acid group, hydroxyl group, alkoxyl group, halogen, nitro group, carbonyl group, amino group, etc. These may be linear or branched.\nThe above alkyl and substituted alkyl are not particularly limited, and include, for example, methyl, ethyl, n-propyl, i-propyl, n-butyl, i-butyl, etc. The above aryl and substituted aryl include a phenyl group, tolyl group, xylyl group, etc. The phenyl group is preferable from the viewpoint of the antifouling property.\nIn the above formula (1), Y represents a primary to tertiary amine or a heterocyclic compound containing a substituted-or unsubstituted nitrogen atom.\nThe above primary to tertiary amine is not particularly limited, and includes, for example, alkyl group-containing amines such as octadecyl amine, capryl amine, lauryl amine, myristyl amine, stearyl amine, dimethyl amine, triethyl amine etc.; acyl group-containing amines such as acetyl amine etc.; alkoxyl group-containing amines such as 3-ethoxypropyl amine, ethoxymethyl amine etc.; aryl group-containing amines such as aniline, hexyl aniline etc.; allyl amine, vinyl amine, aminostyrene, etc.\nThe above heterocyclic compound containing a nitrogen atom may be polycyclic where one ring is preferably a 4- to 10-memberred ring. The heterocyclic compound can also contain an oxygen atom and/or a sulfur atom as the hetero atom in addition to the nitrogen atom thereof, and the number of hetero atoms including the nitrogen atom may be 1 to 4.\nIn the relationship with the above primary to tertiary amine, the heterocyclic compound containing a nitrogen atom refers to the one having the nitrogen atom as a hetero atom constituting the heterocyclic ring, and encompasses both saturated rings and condensed rings containing a double bond in its ring such as pyridine etc.\nThe above heterocyclic compound containing a nitrogen atom is not particularly limited, and includes, for example, pyridine, pyrrole, pyrazol, isoxazole, pyrazine, pyrimidine, imidazole, quinoline, indole, carbazole, phenanthroline; and heterocyclic compounds having a substituent group, such as ethyl pyridine, nicotinic acid, octylnicotinate, vinyl pyridine, allyl pyridine, etc. The above substituent group includes those specifically mentioned hereinabove.\nThe number of carbon atoms in the Y as a whole is preferably 1 to 30. If the number of carbon atoms exceeds 30, the antifouling property will be deteriorated when the composition is incorporated into the antifouling coating so that it will be necessary to be incorporated in a large amount, in some cases. Further, the resulting composition may be hardly dissolved in an organic solvent, and is inferior in handling properties. The number of carbon atoms therein is more preferably 1 to 20.\nIn the present invention, the above Y is preferably a primary amine, pyridine, or pyridine having a substituted group.\nThe boron-containing polymer having at least one group represented by the above formula (1) has preferably a number average molecular weight of 2000 to 40000. If the molecular weight is less than 2000, the resultant resin composition is poor in the film forming ability and may cause abnormalities in the coating film, such as cracking, peeling etc., while if the molecular weight exceeds 40000, the viscosity of the coating will be increased to cause inferior in handling properties, in some cases.\nThe number average molecular weight mentioned above is determined by gel permeation chromatography using polystyrene standards.\nThe boron-containing polymer having at least one group represented by the above formula (1) can be produced by reacting a polymer having at least one acid group in a side chain or at a terminal thereof, with the nitrogen atom-coordinated boron compound represented by the above formula (3). This production method also falls under the scope of the present invention.\nThe above acid group is preferably a carboxylic acid group, a sulfonic acid group or a phosphoric acid group. It is more preferably a carboxylic acid group.\nThe polymer having at least one acid group in a side chain or at a terminal thereof is not particularly limited insofar as it has at least one acid group described above, and examples thereof include an acrylic resin, polyester resin, alkyd resin, epoxy resin, etc. The acrylic resin is preferable from the viewpoint of the antifouling property.\nThe production method of the polymer having at least one acid group in a side chain or at a terminal thereof can be a process known in the art. For example, the acrylic resin can be obtained by copolymerizing a polymerizable monomer having an acid group with other monomers as necessary.\nThe above polymerizable monomer having an acid group includes monomers having a carboxylic acid group, such as methacrylic acid, acrylic acid, itaconic acid, monoalkyl itaconate, monoalkyl maleate (alkyl in the monoalkyl residue includes, for example, C1-20 alkyl such as methyl, ethyl, butyl, 2-ethylhexyl, etc.), half esters of an OH group-containing polymerizable unsaturated monomer and an acid anhydride, for example, half esters of 2-hydroxyethyl (meth)acrylate and succinic anhydride, maleic anhydride, or phthalic anhydride; monomers having a sulfonic acid group, such asp-styrene sulfonic acid, 2-methyl-2-acrylamide propane sulfonic acid etc.; and monomers having a phosphoric acid group, such as phosphoxypropyl methacrylate, 3-chloro-2-phosphoxypropyl methacrylate, phosphoxyethyl methacrylate etc.\nThe above-mentioned other monomers are not particularly limited and include, for example, (meth)acrylates containing 1 to 30 carbon atoms preferably 1 to 20 carbon atoms in the ester moiety thereof, such as methyl (meth)acrylate, ethyl (meth)acrylate, i-propyl (meth)acrylate, n-butyl (meth)acrylate, i-butyl (meth) acrylate, t-butyl (meth) acrylate, 2-ethylhexyl (meth)acrylate, lauryl (meth)acrylate, stearyl (meth)acrylate, 2-hydroxypropyl (meth)acrylate, 2-hydroxyethyl (meth)acrylate, phenyl (meth)acrylate, cyclohexyl (meth)acrylate, polyethylene glycol mono (meth) acrylate, etc.; and (meth)acrylamide, styrene, vinyl acetate, acrylonitrile, etc.\nThese monomers may be used singly or in combination of two or more species.\nThe above copolymerization method is not particularly limited, and for example, mention can be made of a method in which the polymerizable monomer having an acid group and other monomers as necessary are mixed with a polymerization initiator to prepare a mixed solution and then the mixed solution is added dropwise to an organic solvent such as xylene, n-butanol etc., and is allowed to react under heating.\nThe acid value of the polymer having at least one acid group in a side chain or at a terminal thereof obtained in the manner described above is preferably 20 to 400 mg KOH/g. Given an acid value of less than 20 mg KOH/g, the antifouling property of the resulting coating film may be insufficient due to a decreased amount of the nitrogen atom-coordinated boron compound bound to the polymer, while given an acid value of more than 400 mg KOH/g, the speed of elution of the resulting coating film in water may be too rapid to secure the antifouling property for a long period of time. The acid value is more preferably 30 to 250 mg KOH/g. The above acid value is a design value and can be determined according to the species and amount of the monomer to be used.\nThe nitrogen atom-coordinated boron compound is represented by the above formula (3). In the above formula (3), R2 groups are the same or different and each represents a substituted or unsubstituted alkyl or aryl, and Y represents a primary to tertiary amine or a heterocyclic compound containing a substituted or unsubstituted nitrogen atom. The above R2 and Y include those specifically mentioned hereinabove.\nThe above nitrogen atom-coordinated boron compound can be produced in high yield as a stable compound having a nitrogen atom coordinated to boron, by reacting a tri-substituted boron compound to which a basic compound such as sodium hydroxide has been added, with the primary to tertiary amine or the heterocyclic compound containing a substituted or unsubstituted nitrogen atom represented by Y. However, a tri-substituted boron compound to which a basic compound has not been added, for example triphenyl boron etc., is often unstable as it is, and therefore when such a boron compound is reacted with the primary to tertiary amine or the heterocyclic compound containing a substituted or unsubstituted nitrogen atom, the yield becomes low.\nThe tri-substituted boron-compound used as a starting material of the above nitrogen atom-coordinated boron compound may be not only a commercially available product, but also a compound obtainable by reacting BF3.O(Et)2 (Et is ethyl) or the like with a Grignard reagent in THF or dry toluene.\nThe reaction of the above polymer having at least one acid group in a side chain or at a terminal thereof with the above nitrogen atom-coordinated boron compound can be conducted by charging the polymer having at least one acid group in a side chain or at a terminal thereof and the nitrogen atom-coordinated boron compound and then heating the mixture at about 50 to 120xc2x0 C., preferably about 80xc2x0 C. The order of the charging mentioned above is not particularly limited, but for example the polymer having at least one acid group in a side chain or at a terminal thereof maybe added dropwise to an organic solvent such as xylene in which the nitrogen atom-coordinated boron compound has been dissolved.\nIn the above reaction, the mixing ratio between the polymer having at least one acid group in a side chain or at a terminal thereof and the nitrogen atom-coordinated boron compound is generally selected such that 0.3 to 3 molecules, preferably 0.7 to 2 molecules, of the above nitrogen atom-coordinated boron compound are added per acid group into which the nitrogen atom-coordinated boron compound is to be introduced.\nAs the production method of a boron-containing polymer having at least one group represented by the above formula (1) in the present invention, there can also be mentioned a production method comprising a step (I) of reacting a polymerizable organic acid with the nitrogen atom-coordinated boron compound represented by the above formula (3) and a step (II) of polymerizing a monomer component comprising a boron-containing polymerizable monomer obtained in the above step (I). This production method also falls under the scope of the invention.\nThe above polymerizable organic acid is preferably a polymerizable organic carboxylic acid, a polymerizable organic sulfonic-acid or a polymerizable organic phosphoric acid. The above polymerizable organic acid is not particularly limited insofar as it is an organic acid having a polymerizable unsaturated double bond, and for example, mention can be made of those specifically mentioned herein above as the polymerizable unsaturated organic acid having an acid group.\nThe nitrogen atom-coordinated boron compound represented by the above formula (3), the reaction conditions in the step (I) and the polymerization conditions in the step (II) also include those specifically mentioned above. Further, those specifically mentioned above as other monomers can also be copolymerized as necessary in the step (II), together with the boron-containing polymerizable monomer obtained in the step (I) A resin composition comprising the boron-containing polymer that can be produced by the production method described above also falls under the scope of the present invention. Further, when a product having the same structure as that of the boron-containing polymer obtained by the production method of the present invention is obtained by a method other than the production method described above, such a product also falls under the scope of the present invention.\nThe resin composition of the present invention may make use of the boron-containing polymer in combination with other binder resins in order to adjust the hydrolyzability and to adjust the physical properties of the coating film. The above other binder resins include, for example, chlorinated paraffin, polyvinyl ether, polypropylene sebacate, partially hydrogenated terphenyl, polyvinyl acetate, poly(alkyl (meth) acrylate), polyether polyol, alkyd resin, polyester resin, polyvinyl chloride, siliconoil, wax, vaseline, liquidparaffin, rosin, hydrogenated rosin, naphthenic acid, fatty acids and divalent metal salts thereof.\nThe incorporated amount of the above-mentioned other binder resins is preferably [other binder resins]:[boron-containing polymer]=0:100 to 70:30 in a mass ratio. A proportion-of the above other binder-resins higher than the above range is not preferable because both excellent long-term antifouling property and anti-cracking property of the coating film will not be secured simultaneously.\nThe resin composition of the present invention can be incorporated suitably into an antifouling coating. The antifouling coating into which the above resin composition is incorporated also falls under the scope of the invention. The incorporated amount of the above resin composition is generally 0.1 to 80% by mass of the solid content in the antifouling coating. If the amount is less than 0.1% by mass, the resulting coating film is inferior in the antifouling property, while the amount is higher than 80% by mass, there may occur inconveniences such as cracking in the coating film. Preferably, the amount is 1 to 60% by mass.\nThe antifouling coating of the present invention may be obtained by adding conventional additives such as an antifouling agent, a plasticizer, a pigment, a solvent etc. to the resin composition described above.\nThe above antifouling agent is not particularly limited and may be any conventional ones, and examples thereof include powders or flakes of metals such as copper, zinc, nickel etc.; metal oxides, metal hydroxides and metal halides such as those of copper, zinc etc.; other metal salts such as cuprous oxide, copper rhodanide etc.; metal carboxylates such as copper naphthenate, copper stearate etc.; metal dithiocarbamates such as zinc dimethyl dithiocarbamate, bis-dimethyl dithiocarbamoyl zinc, ethylene bisdithiocarbamate etc.; thiuramdisulfides such as tetramethyl thiuram disulfide etc.; sulfamides such as phthalylsulfathiazole, sulfur ethidole, sulfanilidepyridine, sulfomethoxine, N,Nxe2x80x2-dimethyl-Nxe2x80x2-phenyl-N-fluorodichloromethyl thiosulfamide etc.; pyrroles and imidazoles, such as glyodine, fentizole, polycide etc.; thioxanes and thixanthons such as terazole, asterol, miron etc.; imides and amides such as nicarbazin, 3,4,5-tribromosalicyl anilide, N-trichloromethyl mercaptophthalimide, 3,5-dinitrobenzamide, 2,4,6-trichloromaleimide, N-fluorodichloromethyl thiophthalimide etc.; sulfur-containing organic compounds and halogen-containing organic compounds such as 2-methylthio-4-t-butylamino-6-cyclopropylamino-s-triazine, 2,4,5,6-tetrachlorophthalonitrile, N,Nxe2x80x2-dimethyl dichlorophenyl urea, 4,5-dichloro-2-n-octyl-3-(2H)isothiazoline, 2-pyridinethiol-1-oxide zinc salt, 2-pyridinethiol-1-oxide copper salt, 2,3,5,6-tetrachloro-4-methylsulfonyl pyridine, 3-iodo-2-propylbutyl carbamate, diiodomethyl p-tolyl sulfone etc.; pyridine-triphenyl borane, stearyl amine-triphenyl borane, and other agrochemicals, medicines, bactericides, etc.\nThe above plasticizer is not particularly limited, and includes, for example, phthalate-based plasticizers such as dioctyl phthalate, dimethyl phthalate, dicyclohexyl phthalate etc.; aliphatic dibasic ester-based plasticizers such as diisobutyl adipate, dibutyl sebacate etc.; glycol ester-based plasticizers such as diethylene glycol dibenzoate, pentaerythritol alkyl ester etc.; phosphate-based plasticizers such as tricresyl phosphate, trichloroethyl phosphate etc.; epoxy-based plasticizers such as epoxylated soybean oil, octyl epoxy stearate etc.; organic tin-based plasticizers such as dioctyltin laurate, dibutyltin laurate etc.; and trioctyl trimellitate, triacetylene, etc.\nThe above pigment is not particularly limited, and includes, for example, extender pigments such as precipitated barium sulfate, talc, clay, chalk, silica white, alumina white, bentonite, etc.; and color pigments such as titanium oxide, zirconium oxide, basic lead sulfate, tin oxide, carbon black, graphite, red iron oxide, chromium yellow, phthalocyanine green, phthalocyanine blue, quinacridone, etc.\nThe above solvent is not particularly limited, and includes, for example, hydrocarbons such as toluene, xylene, ethyl benzene, cyclopentane, octane, heptane, cyclohexane and white spirit; ethers such as dioxane, tetrahydrofuran, ethylene glycol monomethyl ether, ethylene glycol monoethyl ether, ethylene glycol dimethyl ether, ethylene glycolmonobutyl ether, ethylene glycol dibutyl ether, diethylene glycol monomethyl ether and diethylene glycol monoethyl ether; esters such as butyl acetate, propyl acetate, benzyl acetate, ethylene glycol monomethyl ether acetate and ethylene glycol monoethyl ether acetate; ketones such as ethyl isobutyl ketone and methyl isobutyl ketone; and alcohols such as n-butanol and propyl alcohol.\nBesides those described above, other additives are not particularly limited, and include, for example, rosin, organic monobasic acids such as monobutyl phthalate and monoctyl succinate etc., camphor, castor oil, etc.\nThe antifouling coating of the present invention can be prepared for example by adding conventional additives such as other binder resins, an antifouling agent, a plasticizer, a coating-abrasion regulator, a pigment, a solvent etc. to the above resin composition comprising the boron-containing polymer according to the present invention and then mixing them by a mixer such as a ball mill, a pebble mill, a roll mill, a sand grind mill, etc.\nA dry coating film can be formed by applying the antifouling coating described above in a usual manner onto the surface of a substrate to be coated and then removing the solvent through evaporation at ordinary temperature or under heating.\nThe resin composition comprising the boron-containing polymer according to the present invention is a novel compound expressing an excellent antifouling property as a binder resin in an antifouling coating, and it is a hydrolysable resin having a so-called self-polishing property where the resin is gradually hydrolyzed in water to allow abrasion of the coating film.\nThe antifouling coating incorporated with the resin composition comprising the boron-containing polymer according to the present invention is superior in the antifouling property of the resulting coating film. The-above antifouling coating can be used suitably on ships, structures in water, fishing nets such as farming nets and stationary nets, water conduits, fishing implements, etc.\nThe resin composition comprising a boron-containing polymer according to the present invention, when incorporated into an antifouling coating, is excellent in the antifouling property of the resulting coating film. The above boron-containing polymer can be synthesized simply and easily according to the production method of the present invention. The above antifouling coating can be used suitably on ships, structures in water, fishing nets such as farming nets and stationary nets, water conduits, fishing implements, etc."}
{"text": "1. Field of the Invention\nThis invention relates generally to a totally self-contained hermetically sealed, immersible device capable of assuming alternately, a non-buoyant and buoyant mode and particularly to a dual buoyancy device which, in its minimum volumetric dimension displaces a volume of liquid whose weight is less than that of the device thus causing it to sink; with an increase to maximum volumetric dimension but with unchanged total weight, the device then displaces a volume of liquid whose weight is greater than that of the device, causing it to rise and float in the liquid.\n2. Prior Art\nBuoyant devices suitable for attachment to various objects to cause them to be lifted to the surface of a liquid are well known. They are usually made of buoyant material or have a shell-like structure which makes them buoyant.\nUnder most circumstances, buoyancy of these devices is reduced so that they will sink, and they may be attached to an object. This is usually accomplished by admitting to the interior of the device, sufficient of the liquid in which the device is immersed so that its buoyancy is reduced and it sinks. After attachment to an object, liquid within the device is pumped or otherwise removed restoring its normal buoyancy; which resulting force lifts the attached object to the surface."}
{"text": "Optical fibers have displaced copper-based connectivity in much of the traditional long-haul and metro telecommunication networks for numerous reasons such as large bandwidth capacity, dielectric characteristics and the like. As higher network speeds for communication networks are required the optical fiber will move deeper into the communication networks toward the electronics located in servers and switches that manage traffic on the communication network. As this migration of optical fiber expands deeper into communication networks new challenges will be encountered for making optical connections to electronics for high-speed communications. By way of example, aligning and maintaining optical alignment for a high-density optical connection with integrated circuits presents unresolved needs for the industry."}
{"text": "Field of the Art\nThe invention relates generally to the field of robotics and more particularly to the field of robotic physical data acquisition.\nDiscussion of Prior Art\nIn the art of robotics, it is a common practice to manually control robotic devices or appliances remotely, for example in radio-controlled vehicles, surveillance drones, and similar devices. It is a common practice to collect data using these devices, for example by utilizing a drone-mounted camera to record photographic information, or microphones to record acoustic information. However, in the art there is no way to unify the gathered data, for example to model buildings and their systems using remotely controlled robotic devices or appliances. What is needed, is a way to combine collected data into an accessible format for use, for example, to enable enterprises to plan, schedule and verify preventive maintenance, supply chain, subcontractors, scheduling, enhanced object model libraries, ongoing management optimization, anomaly detection, and other such purposes. For example, HVAC motor sensors could be timed to optimize scheduling intervals."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical recording/reproducing method and apparatus for recording information onto and reproducing information from a phase-change storage medium, such as a DVD-rewritable.\n2. Description of the Related Art\nDVDs (digital versatile disks) are read-only optical recording media and have their rewritable equivalents, DVD-rewritable disks. DVD-rewritable disks are phase-change storage media that use a process called the phase change technology to write and erase data.\nThe process writes data to a DVD-rewritable disk by focusing a high-intensity laser beam on a layer of a phase-change material embedded in a plastic base or a substrate of the disk. In its original state of the phase-change material, this layer has a rigid crystalline structure. The laser beam selectively heats areas of a surface portion of the disk to a high temperature. Where the beam strikes, the heat melts the crystals to a non-crystalline, or amorphous phase. These areas reflect less light than the unchanged area surrounding them.\nLater, when a weaker laser beam, used only to read data from the recorded disk, strikes the non-crystalline area, the beam is scattered and hard to be picked up by the light-sensitive diode in the read head of the disk drive. With the lower reflectance, these areas become marks, representing xe2x80x9c1xe2x80x9ds. Areas that are not heated are more reflective areas, representing xe2x80x9c0xe2x80x9ds. When the read laser beam strikes the areas, it is reflected directly to the light-sensitive diode of the read head, creating an electrical current that is sent to the controller in the drive. The controller interprets the pattern of electrical pulses, decodes the data that they represent, checks the data for error, and sends the data to a computer.\nTo erase data or to change a mark back to crystalline phase, DVD-rewritable disk drives use a lower-energy laser beam to heat marked areas to a comparatively low temperature. This amount of heat is below the melting point of the phase-change material, but it still loosens up the phase change media so that it can recrystallize to its original state.\nRecently, there is the demand for using a DVD player or the like to reproduce data from a DVD-rewritable disk (or a phase-change storage medium). The DVD player is originally designed to use a laser beam to read the pits of DVD or DVD-ROM disks (or the read-only optical storage media). On the read-only optical storage media, various data blocks are recorded in a seamless formation. If data blocks are recorded on the phase-change medium in a seamless formation that is essentially the same as that of the read-only optical storage media, data can be easily reproduced from the phase-change medium by using the DVD player without a need for substantial changes of the DVD player. This will increase the operability and compatibility of DVD-rewritable media, which considerably benefits the users.\nGenerally, in optical disk drives using the phase-change media, if data is repeatedly written to a same data block on the phase-change disk, the recording properties of the disk will deteriorate and a defect in the data block will take place. In such a case, the optical disk drives do not effectively reproduce the data from the disk. There are three types of the deterioration pattern that may be revealed in the phase-change media.\nThe first type is that, if data is repeatedly written to a head end of one data block, defective areas are transferred from the head end of the data block to the following areas in the forward direction. The second type is that, if data is repeatedly written to a tail end of one data block, defective areas are transferred from the tail end of the data block to the preceding areas in the backward direction. The final type is that, if a same pattern of marks is repeatedly recorded onto a same position of one data block, defective areas are transferred from that position in both the forward direction and the backward direction.\nThe larger the energy of the laser beam applied to overwrite recording marks at the same position of the phase-change disk, the severer the deterioration of the recording properties of the disk.\nIn order to eliminate the above problem, a conventional optical recording/reproducing apparatus is configured to vary the position of each complete data block in the disk when overwriting data. See, for example, Japanese Patent No.2720876 (which corresponds to Japanese Laid-Open Patent Application No.9-180191). However, the changes of the position of each complete data block is unsuitable for an optical disk drive that records data blocks onto the optical disk in a seamless formation. In the conventional apparatus, the sync positions of data blocks are greatly shifted from each other and it is difficult to create a read clock signal the phase of which is accurately synchronized with the read clock signal of the following data block when reproducing data from the optical disk.\nIn order to eliminate the above problem, another conventional optical recording/reproducing apparatus is configured to add dummy data to the head end and tail end of each data block and vary the position of the entire data block in the disk so as to protect the recorded data. See, for example, Japanese Laid-Open Patent Application No. 9-219022. However, the changes of the position of the entire data block are unsuitable for an optical disk drive that records data blocks onto the optical disk in a seamless formation.\nAnother conventional optical recording/reproducing apparatus is configured to produce a write clock signal precisely synchronized with rotation of an optical disk. See, for example, Japanese Laid-Open Patent Application No.10-293926. However, the write clock signal produced is unsuitable for an optical disk drive that records data blocks onto the optical disk in a seamless formation. It is difficult for the above conventional apparatus to create a read clock signal the phase of which is accurately synchronized with the read clock signal of the following data block when reproducing data from the optical disk.\nIn addition, in the above conventional apparatus, when reproducing data from the optical disk, there is a problem in that the phase of the read clock signal obtained from the preceding data block and the phase of the read clock signal obtained from the next data block are apart from each other by nearly 180 degrees. It is highly possible that the data recorded at the boundaries between the adjacent data blocks on the optical disk are lost.\nAs described above, in the above-mentioned conventional apparatuses, the gaps between adjacent data blocks are provided, or the dummy data are added to the head end and tail end of each of the data blocks. It is difficult for the above-mentioned conventional apparatuses to record data blocks onto the phase-change disk in a seamless formation while effectively reducing the deterioration of the recording properties of the disk.\nIn order to overcome the problems described above, preferred embodiments of the present invention provide an improved optical recording/reproducing method and apparatus that accurately records data blocks onto the phase-change disk in a seamless formation while effectively reducing the deterioration of the recording properties of the disk.\nAccording to one preferred embodiment of the present invention, an optical recording/reproducing method for accessing an optical storage medium in which binary data are recorded by applying a light beam to a recording layer of the storage medium and changing a state of the recording layer on a disk drive, includes the steps of: receiving a data signal to record a sequence of data blocks onto the storage medium, each data block having a predetermined block length from one end of the data block to the other end; varying each of the data blocks in length such that data bits are arbitrarily added to or removed from at least one of the ends of the data block in order to allow the data blocks to be written to the storage medium in a seamless formation, the data bits having a predetermined length; and recording the resulting data blocks on the storage medium at fixed positions.\nIn the optical recording/reproducing method of the above preferred embodiment, the data blocks are recorded onto the optical disk in the seamless formation, and the optical recording/reproducing method is effective in accurately recording data blocks onto the phase-change disk in a seamless formation and in reducing the deterioration of the recording properties of the disk.\nAccording to another preferred embodiment of the present invention, an optical recording/reproducing method for accessing an optical storage medium includes the steps of: receiving a data signal to record a sequence of data blocks onto the storage medium, each data block on the storage medium having a predetermined block length from one end of the data block to the other end and being divided into a number of frames each of which has an information unit that always appears at a same position within the frame with a fixed length when recording of data onto a same portion of the storage medium is repeated; and forming, when marks are recorded onto the storage medium, the information unit of each of the frames in the sequence of data blocks such that the information unit consists of marks of a plurality of sub-units.\nThe optical recording/reproducing method of the above preferred embodiment is effective in reducing the deterioration of the recording properties of the optical disk on which data blocks are recorded in a seamless formation. It is possible to increase the recording repeatability and the operational life of the optical disk.\nAccording to another preferred embodiment of the present invention, an optical recording/reproducing method for accessing an optical storage medium in which binary data are recorded by applying a light beam to a recording layer of the storage medium and changing a state of the recording layer on a disk drive, includes the steps of: calibrating the disk drive with the storage medium so that a phase error of a readout signal relative to a phase of a recording pulse based on a position-information carrying signal obtained from the storage medium is detected; and adjusting the phase of the recording pulse so as to reduce the phase error of the readout signal.\nIn the optical recording/reproducing method of the above preferred embodiment, the read clock signal and the write clock signal can accurately accord in phase, and it is possible to prevent the occurrence of a read error when reproducing data from the disk.\nAccording to another preferred embodiment of the present invention, an optical recording/reproducing apparatus for accessing an optical storage medium includes: a read clock generating unit for generating a read clock signal in response to recorded data on the storage medium; a write clock generating unit for generating a write clock signal in response to a position-information carrying signal detected from the storage medium; a detection unit for detecting a phase difference between the read clock signal and the write clock signal when writing data blocks to the storage medium in a seamless formation; and a determination unit for determining whether a phase of the read clock signal is delayed or advanced based on the phase difference detected by the detection unit, wherein the determination unit produces a control signal to adjust a phase of the position-information carrying signal.\nIn the optical recording/reproducing apparatus of the above preferred embodiment, when the optical disk is recorded in a seamless formation, the read clock signal and the write clock signal can accurately accord in phase, and it is possible to prevent the occurrence of a read error when reproducing data from the disk."}
{"text": "This invention relates to the field of stimulated x-ray emission and, more particularly, to the field of x-ray emission caused by a lasing type of action.\nThere is a recognized need for a coherent x-ray generator, but the short life of the excited state of atoms which emit x-rays has presented a problem difficult to overcome. Recently, there have been developed techniques for overcoming the effects of very short spontaneous emission lifetimes by causing a population inversion to sweep along the lasing direction at the speed of light. In the copending U.S. application Ser. No. 393,293 of W. B. McKnight et al., filed Aug. 30, 1973 and entitled \"X-Ray Laser,\" there is disclosed a system wherein an ion beam of completely stripped nuclei, such as He.sup.++, Li.sup.+++, etc., is passed through a foil so as to effect resonant charge pickup in an excited state. The stripped nuclei are swept across the foil target at the speed of light. When the beam strikes the foil target, a population inversion in the states between which laser action occurs is established by means of beam foil excitation. The foil is located in one of the two plates which are aligned parallel to the incident direction of the beam. A voltage pulse source is connected to the end of the two plates so as to cause a voltage pulse to travel down the plates at the speed of light. The resultant electric field causes deflection of the beam onto the negative plate where the thin foil is located, so that the beam sweeps along and into the foil at the same speed as the velocity of propagation of the wavefront of stimulated emission. In this manner, a volume element with a large density of excited states on the exit side of the foil is produced due to the resonant charge exchange between the beam and the hydrogen in the foil. Since the beam sweeps downstream in coincidence with the stimulated emission wavefront, the traveling stimulated emission constitutues an x-ray laser action.\nWhile the described technique is considered promising, there are problems associated with the production and use of the thin hydrogen-rich foil strip. The foil thickness needed to maximize the desired population inversion is so thin that such foils are difficult to make which are still thick enough so that multiple scattering will give a suitable Doppler width. In addition, the lifetime of the foils are short since they are destroyed by the impinging ion beam.\nIt is an object of this invention to overcome the disadvantages set forth."}
{"text": "Throughout this application various publications are referred to by partial citations within parenthesis. Full citations for these publications may be found at the end of the specification immediately preceding the claims. The disclosures of these publications, in their entireties, are hereby incorporated by reference into this application in order to more fully describe the state of the art to which this invention pertains.\nMolecular cloning efforts have provided primary amino acid sequence and signal transduction data for a large collection of serotonin receptor subtypes. These include five cloned 5-HT.sub.1 -like receptors, three cloned 5-HT.sub.2 receptors, and one 5-HT.sub.3 receptor. The 5-HT.sub.1 subfamily includes: 5-HT.sub.1A (Fargin, 1988; Kobilka, 1989), 5-HT.sub.1B /5-HT.sub.1D.beta. (Weinshank et al., 1991; Demchyshyn et al., 1992; Jin et al., 1992; Adham et al., 1992; Maroteaux et al., 1992; Voight et al., 1991), 5-HT.sub.1D.alpha. (Branchek et al. 1991; Hamblin and Metcalf, 1991; Weinshank et al., 1992), 5-HT.sub.1E (Levy et al., 1992; McAllister et al., 1992; Zgombick et al., 1992) and 5-HT.sub.1F (Adham et al., 1993). All five have been shown to couple to the inhibition of adenylate cyclase activity. The 5-HT.sub.2 family includes the 5-HT.sub.2 receptor (Pritchett et al., 1988), 5-HT.sub.1C (Julius et al., 1989) and 5-HT.sub.2F (Rat Stomach Fundus; Foquet et al., 1992; Kursar et al., 1992). These receptors all couple to phosphoinositide hydrolysis. The 5-HT.sub.3 receptor is a ligand-gated ion channel (Maricq et al., 1991).\nAlthough this work represents enormous success, the absence of molecular biological information on the 5-HT.sub.4 receptors, which have been shown in native tissues to couple to the activation of adenylate cyclase as a primary mode of signal transduction (Dumius et al., 1988; Bockaert et al., 1990), is apparent. In a previous copending application (U.S. Ser. No., 971,690, filed Nov. 3, 1992), we reported the cloning of the first mammalian 5-HT receptor that couples to the stimulation of adenylate cyclase activity which we named 5-HT.sub.4B. The pharmacological properties of this receptor indicated that it was similar to a series of functionally defined 5-HT receptors described in the porcine vena cava (Trevethick et al., 1984), cat saphenous vein, coronary arteries (Cushing and Cohen, 1992), and several vascular dilatory effects (Mylecharane and Phillips, 1989). However, the classically defined 5-HT.sub.4 receptor remained to be cloned. We now report the cloning of the pharmacologically-defined 5-HT.sub.4 receptor which we have called 5-HT.sub.4A. This receptor also stimulates adenylate cyclase activity but unlike 5-HT.sub.4B, is sensitive to a series of benzamide derivatives which act as agonists or partial agonists at this subtype. The presence of this subtype in the brain, particularly in areas such as the hippocampus, indicates a potential role in cognitive enhancement. In addition, the 5-HT.sub.4 receptor has been described functionally in the heart, adrenal, bladder, and alimentary canal indicating potential roles in achalasia, hiatal hernia, esophageal spasm, irritable bowel disease, postoperative ileus, diabetic gastroparesis, emesis and other diseases of the gastrointestinal tract, as well as in cardiac, urinary, and endocrine function."}
{"text": "The present invention relates to a magneto-rheological (xe2x80x9cMRxe2x80x9d) fluid damper, and more particularly, to a linearly-acting MR fluid damper suitable for vibration damping in a vehicle suspension system.\nMR fluids are materials that respond to an applied magnetic field with a change in Theological behavior (i.e., change in formation and material flow characteristics). The flow characteristics of these non-Newtonian MR fluids change several orders of magnitude within milliseconds when subjected to a suitable magnetic field. In particular, magnetic particles noncolloidally suspended in fluid align in chain-like structures parallel to the applied magnetic field, changing the shear stress on adjacent shear surfaces.\nDevices such as controllable dampers benefit from the controllable shear stress of MR fluid. For example, linearly-acting MR fluid dampers are used in vehicle suspension systems as vibration dampers. At low levels of vehicle vibration, the MR fluid damper lightly damps the vibration, providing a more comfortable ride, by applying a low magnetic field or no magnetic field at all to the MR fluid. At high levels of vehicle vibration, the amount of damping can be selectively increased by applying a stronger magnetic field. The controllable damper lends itself to integration in vehicle suspension systems that respond to vehicle load, road surface condition, and driver preference by adjusting the suspension performance.\nIn some applications, linearly-acting MR fluid dampers use a piston assembly that moves within a damper body tube having a cylindrical reservoir that separates a volume of MR fluid into a compression chamber and an extension chamber. The piston assembly has a piston core positioned within a flux ring to form an annular flow gap therebetween. Relative motion between the damper body tube and the piston assembly is dampened by a flow of the MR fluid through the flow gap from one chamber to another caused by the relative motion.\nAlignment of the flux ring is critical for optimum performance. Ideally, the piston assembly should move freely within the reservoir in the damper body tube without friction or binding. In addition, the radial width and concentricity of the annular flow passage must be precisely set and maintained along the axial length of the passage throughout the operation to ensure optimum, predictable control of the damping. Consequently, the flux ring must be correctly aligned with the piston core.\nAttachment elements have been suggested to provide flux ring alignment with nonmagnetic bridge elements. In particular, perforated end plates are aligned above and below the flux ring and piston core. These attachment elements have several potential problems. First, the attachment elements increase the length of the piston assembly. Consequently, less travel distance is available for the piston to move within the cylindrical reservoir of the damper body tube. Second, the attachment elements require tight manufacturing tolerances in order to correctly align the flux ring to the piston core. Third, such attachment elements often include tabs or other projections that increase the drag as the piston moves, which may be undesirable. Fourth, the attachment elements have numerous components and require manufacturing operations such as spot welding. Therefore, such attachment elements are costly to manufacture and time consuming to assemble.\nConsequently, there is a need for an improved piston assembly suitable for use in a magneto-rheological (MR) fluid damper.\nThe present invention addresses the above need by providing an improved piston assembly for a linearly-acting MR fluid damper. The piston assembly of the present invention confines a flux ring within the functional length of the piston assembly without significantly restricting fluid flow, thereby providing optimum performance with minimal piston length. Further, the part count of the piston assembly is reduced; and the piston assembly is easier to assemble in a desired alignment. Thus, the piston assembly of the present invention is of a simpler construction than known damper pistons that can be manufactured for less cost.\nAccording to the principles of the present invention and in accordance with the described embodiment, the present invention provides a piston assembly for use with an MR fluid damper. The piston assembly has a flux ring positioned in a desired alignment with a piston core to form an annular flow gap between the flux ring and the piston core. The piston core is secured to the flux ring in the desired alignment by a plurality of projections extending across the flow gap between an inner surface of the flux ring and an outer surface of the piston core. Thus, the flux ring is secured on the piston ring without using expensive, high precision attachment components; and the piston assembly is able to utilize its full length, thereby providing optimum performance with a minimum of length piston.\nIn one aspect of the present invention, the projections are molded through attachment passages in the flux ring. In a further aspect of the invention, the attachment passages are holes intersecting the inner and outer surfaces of the flux ring.\nIn another embodiment of the invention, a method is provided for making a piston assembly for use with an MR fluid damper. The method comprises first, fixing a piston core and a flux ring in a desired alignment forming a flow gap therebetween; and then, forming a plurality of projections in the flow gap between the flux ring and the piston core to secure the flux ring on the piston core in a desired alignment."}
{"text": "1. Field of the Invention\nExample embodiments according to the present general inventive concepts relate to semiconductor memory devices. More particularly, the example embodiments relate to methods of merging blocks in semiconductor memory devices, and semiconductor memory devices that perform methods of merging blocks.\n2. Description of the Related Art\nSemiconductor memory devices can be roughly divided into volatile memory devices and nonvolatile memory devices according to whether they retain stored data when disconnected from power. Recently, as electronic devices require lower power consumption and smaller size, a NAND flash memory device has been widely used as the nonvolatile memory device. In the NAND flash memory device, a program operation is performed on a page basis, and an erase operation is performed on a block basis.\nGenerally, since the NAND flash memory device, unlike hard disc drives, does not support an overwrite operation, data modification in the NAND flash memory device should be performed in a manner that programs data of a page to be modified into a free page of the same block or another block with data modification, and assigns an address of the newly programmed page to the modified data by a flash translation layer (FTL). Accordingly, the page that previously stored the data becomes an invalid page, and the page that stores the modified data becomes a valid page. That is, the invalid page may coexist with the valid page in one block.\nTo obtain free blocks, a block merge operation can be performed by merging valid pages distributed in a plurality of blocks and by erasing unnecessary invalid pages. However, overall performance of the NAND flash memory device may drastically decrease due to the time required to perform the block merge operation."}
{"text": "This invention relates to a new foam product suitable for use as insulation in the building industry. Involved are resin systems which have been previously investigated but not used together in the system disclosed herein. One of these systems is based upon the urethane formation reaction wherein an isocyanate is reacted with a compound containing hydroxyl groups. A second reaction involved is the trimerization of isocyanates, a reaction used in the production of polyisocyanurate products. Finally, a third reaction involved is the curing of an unsaturated polyester resin with a cross-linking vinyl monomer in the presence of a free radical catalyst. The present invention combines these reactions in a manner not previously utilized by mixing two or more reactant streams in the manner to be fully described hereinafter.\nOne example of such prior operation is Graham et al., U.S. Pat. No. 3,860,537 (1975) where a stepwise reaction is used. In the system disclosed in that patent, a mixture of polyurethane precursors comprising a polyol and a polyisocyanate is mixed with polyester resin precursors comprising an unsaturated polyester and an unsaturated monomer. These mixtures of precursors are mixed and partially cured so that the precursors in the mixture react such that the mixture of the polyfunctional compound and the ethylenically unsaturated polyester are substantially reacted with the polyisocyanate, and the ethylenically unsaturated polyester and the unsaturated monomer are substantially uncopolymerized. The patent describes subsequently carrying out the second polymerization. In the system disclosed herein, all of the reactions take place at substantially the same time.\nAdditional work in the field is represented by Hutchinsen et al., U.S. Pat. No. 3,886,229 (1975) although this patent does not disclose foam products. In the system of Hutchinson et al., as in Graham et al., a polymeric product is formed by mixing precursors of a polyurethane comprising at least one diisocyanate and at least one difunctional compound with a second component comprising precursors of a cured polyester resin comprising an unsaturated polyester and at least one monomer copolymerizable with the unsaturated groups in the polyester. It will be apparent from the description of the invention disclosed herein that different preliminary mixtures are made prior to the mixing and reaction of the various ingredients.\nIn view of the above, it is an object of this invention to provide a new system for the production of a new insulating foam board.\nA further object of this invention is to provide an insulating foam board based on an unsaturated polyester resin, which is capable of high filler loadings.\nA further object of the invention is to provide an insulation foam board having a flame spread less than 50, preferably less than 25, and smoke generation less than 450 in the ASTM E84 test.\nOther objects and advantages of the present invention will be apparent to those skilled in the art upon reading this disclosure."}
{"text": "Conventionally, photo organization applications provide databases that store information associated with a collection of photos. The photo organization applications provide metadata that may be used by a user to organize the collections of photos. The photo organization application conventionally utilizes a database to associate the metadata with the multimedia files. Because the metadata is not written back to the file, the user is limited to using the photo organization application on the client that created the metadata.\nSome photo organizations application may allow the user to expressly indicate that metadata should be stored with the photos and the database. When the metadata is stored with the photos, a user may transfer the photos to a different client and retain the organization structure created by the photo organization application. However, the photo organizations applications that allow the metadata to be written back to photos suffer from long latency periods during which the client is unable to access the photos. Furthermore, the photo organization applications do not provide a user with the status of the metadata write backs. Consequently, a client may exit the photo organization application prior to the completing the write back operations for the photos. The premature exit of the photo organization application may cause the photo to lose some metadata."}
{"text": "In recent years, with the development of the Internet and communication devices, the numbers of pieces of data and amounts of message data transmitted from personal smart devices are also increasing. In a system that processes these message data, it is necessary to construct a system by using distributed data storage in order to improve extensibility and serviceability of the system.\nA distributed processing type message queue system that stores message data in distributed data storages and processes message data in a distributed manner includes a plurality of message queue servers, for example. Each message queue server includes a storage device that stores message data. The plurality of message queue servers share storage devices and provide a message queue for each storage device between a client and a service providing server.\nAs the plurality of message queue servers share storage devices with one another, large-capacity storage can be implemented. In addition, the storage capacity can be easily increased by adding a message queue server. Also, serviceability can be improved by duplicating message data between message queue servers.\nThere is a related technique disclosed in Patent Literature 1. The PTL 1 discloses a technique for controlling a processing order based on correlation between different messages within a same message queue."}
{"text": "Programmable resistance memory elements formed from materials that can be programmed to exhibit at least a high or low stable ohmic state are known in the art. Such programmable resistance elements may be programmed to a high resistance state to store, for example, a logic ONE data bit or programmed to a low resistance state to store a logic ZERO data bit.\nOne type of material that can be used as the memory material for programmable resistance elements is phase-change material. Phase-change materials may be programmed between a first structural state where the material is generally more amorphous (less ordered) and a second structural state where the material is generally more crystalline (more ordered). The term \"amorphous\", as used herein, refers to a condition which is relatively structurally less ordered or more disordered than a single crystal and has a detectable characteristic, such as high electrical resistivity. The term \"crystalline\", as used herein, refers to a condition which is relatively structurally more ordered than amorphous and has lower electrical resistivity than the amorphous state.\nThe phase-change materials may be programmed between different detectable states of local order across the entire spectrum between completely amorphous and completely crystalline states. That is, the programming of such materials is not required to take place between completely amorphous and completely crystalline states but rather the material can be programmed in incremental steps reflecting (1) changes of local order, or (2) changes in volume of two or more materials having different local order so as to provide a \"gray scale\" represented by a multiplicity of conditions of local order spanning the spectrum between the completely amorphous and the completely crystalline states. For example, phase-change materials may be programmed between different resistive states while in crystalline form.\nA volume of phase-change material may be programmed between a more ordered, low resistance state and a less ordered, high resistance state. A volume of phase-change is capable of being transformed from a high resistance state to a low resistance state in response to the input of a single pulse of energy referred to as a \"set pulse\". The set pulse is sufficient to transform the volume of memory material from the high resistance state to the low resistance state. It is believed that application of a set pulse to the volume of memory material changes the local order of at least a portion of the volume of memory material. Specifically, it is believed that the set pulse is sufficient to change at least a portion of the volume of memory material from a less-ordered amorphous state to a more-ordered crystalline state.\nThe volume of memory material is also capable of being transformed from the low resistance state to the high resistance state in response to the input of a single pulse of energy which is referred to as a \"reset pulse\". The reset pulse is sufficient to transform the volume of memory material from the low resistance state to the high resistance state. While not wishing to be bound by theory, it is believed that application of a reset pulse to the volume of memory material changes the local order of at least a portion of the volume of memory material. Specifically, it is believed that the reset pulse is sufficient to change at least a portion of the volume of memory material from a more-ordered crystalline state to a less-ordered amorphous state.\nThe use of phase-change materials for electronic memory applications is known in the art. Phase-change materials and electrically programmable memory elements formed from such materials are disclosed, for example, in U.S. Pat. Nos. 5,166,758, 5,296,716, 5,414,271, 5,359,205, 5,341,328, 5,536,947, 5,534,712, 5,687,112, and 5,825,046 the disclosures of which are all incorporated by reference herein. Still another example of a phase-change memory element is provided in U.S. patent application Ser. No. 09/276,273, the disclosure of which is also incorporated herein by reference.\nIt is important to be able to accurately read the resistance states of programmable resistance elements which are arranged in a memory array. The present invention describes an apparatus and method for accurately determining the resistance states of programmable resistance elements arranged as memory cells in a memory array. Background art circuitry is provided in U.S. Pat. No. 4,272,833 which describes a reading apparatus based upon the variation in the threshold levels of memory elements, and U.S. Pat. No. 5,883,827 which describes an apparatus using a fixed resistance element to generate reference signals. Both U.S. Pat. No. 4,272,833 and U.S. Pat. No. 5,883,827 are incorporated by reference herein."}
{"text": "Soft drink bottles are ordinarily packaged by bottlers in cases holding several plastic or glass bottles for shipment to retailers or for storage. The cases are customarily stacked on top of the bottles in the lower tray for storage in warehouses. The trays must therefore be particularly stable in order to remain standing in storage. The bottles are generally stored in plastic low depth cases in which the side walls of the cases are lower than the height of the bottles and in which the bottles support the weight of the cases stacked on top of the bottles. When cases of these are stacked for display, the only bottles which can be removed from the stack are those stored in the top case.\nCurrently in order for a bottle to be displayed, it must be removed from the tray and placed in a display rack. If the trays are stacked on the bottles for display, the bottles can only be removed from the top tray. The top tray is removed from the stack when empty to provide access to the bottles on the next tray. Trays stacked on the tops of the bottles in the lower trays are unstable, particularly when bottles are removed from the lower trays while in storage."}
{"text": "Men's underwear typically suffers from several deficiencies, including but not limited to, a propensity for the scrotum to contact and stick to the inner thigh and a high degree of contact between the inner thighs of the wearer which can lead to friction, chaffing, and discomfort while sitting, walking, running, or skating or other numerous athletic activities, a lack of support for the genitalia, limited freedom of movement and natural resting position for the genitalia and a lack of moisture control.\nVarious underwear designs in the prior art have tried to address one or more of these problems. For example, sports boxer-briefs are typically tight and form-fitting resulting in compressing the genitalia. However, they support the scrotum too close to the body leading to any combination of scrotum contact with inner thighs and legs, and moisture problems resulting from perspiration and skin irritation. Typical loose fitting boxer shorts provide ample freedom of movement but suffer from a propensity for the scrotum to contact and stick to the inner thigh, and lack of support and comfort. Further, the prior art typically does not use performance fabrics as those fabrics are better described below. Typical underwear briefs prevent the scrotum from contacting with the legs, but suffer from unnaturally restricting the genitalia into a small compartment pressing with an unnaturally close proximity to the body with the resultant increase in moisture problems and lack of comfort. To date, no underwear product has successfully addressed the problem of a lack of isolation of scrotum from contacting inner thighs and legs while increasing comfort and support attributes for male genitalia.\nIn the prior art applicant is aware of U.S. Pat. No. 5,029,345 which issued to Angheluta, et al, on Jul. 9, 1991, entitled Man's Underwear with an Integral Codpiece-like Construction. Angheluta discloses a man's garment, particularly a man's underwear garment, providing a codpiece compartment secured to a front of a trunk panel of the garment. The trunk panel provides an opening for communicating a man's genitals therethrough into the codpiece compartment. The compartment loosely holds the man's genitals. The trunk panel and a slot 24 provides for forward biasing and separation of the genitals by constricting of the sides of the slot behind the testicles and heat insulation to reduce the thermal influence of the man's torso on the man's reproductive organs. Angheluta describes the use of compartment 40 which is a loose fitting compartment around the genitals, it being taught that such a loose fitting compartment prevents compression of the genitals against the male's lower trunk, such as occurs in shorts known to the art, and that by properly selecting the lengths of the front, right and left side panels 32, 38, and 36 a suitably loose compartment 40 is formed. Angheluta concludes that by reducing compression of the genitals against the trunk, the thermal influence of the trunk to the genitals is reduced, and pressure on the genitals is reduced."}
{"text": "Operations performed in surgery of the spine may concern the cervical region (neck), the thoracic region or, more frequently, the lumbar region.\nWhen there is an instability, for example sliding of one vertebra relative to the adjacent vertebrae, an operation for stabilizing the spinal column may entail implanting metal material in the form of screws connected to each other by bars or plates. These implants form scaffolding that acts as a stabilizer for the spinal column.\nIn more recent techniques, the stabilization of the column can be obtained by means of an intervertebral implant composed of a stabilizing wedge, a flexible link in the form of a textile braid, a movable assembly and a locking member. The stabilizing wedge is intended to be positioned between the spinous processes of two consecutive vertebrae, that is to say adjacent vertebrae, that are to be stabilized. The flexible link (for example a textile braid) encloses the spinous processes. The movable assembly is adapted to come into engagement with the stabilizing wedge in such a way as to immobilize the flexible link with respect to the stabilizing wedge. This blocking is achieved by clamping the flexible link between the movable assembly and the stabilizing wedge. The locking member (for example a screw) is adapted to lock the engagement of the movable assembly with the stabilizing wedge, and thus the final blocking of the flexible link, which results directly from this."}
{"text": "The field of the invention is exercise devices, and the invention relates more particularly to lightweight exercising devices useful for exercising the upper arms, shoulders, calves or the like. Elastic members have been used in exercising devices for many years. An early device is shown in the Bailey U.S. Pat. No. 829,754. This device, however, is not easy to adjust and would be dangerous if the elastic member broke. A spring-type exercising device is shown in the Weiss et al. U.S. Pat. No. 3,815,904 where a plurality of rubber bands are supported between a clamped lower bar and a U-shape upper bar. This device may be used for exercising the arms, back or legs depending on where the lower bar is attached. Like the Bailey device, the Weiss device has no safety provisions if one or more of the rubber bands breaks, and it is, likewise, not easy to adjust. A hand exercising device is shown in the Palmer U.S. Pat. No. 4,783,067. This uses a pair of rubber bands on the exterior of the device and although the bands are very easily changed, the exercising device is quite limited to either exercising the hand or the chest. No safety provisions are provided in the event one of the rubber bands breaks.\nThe United Kingdom Patent No. 2,148,136A shows an arm exercising apparatus which has a hand grip to which one or more elastic cords may be affixed. The other end of the elastic cords are snapped to a support bar. No provisions are made for safety in the event of a breakage of one of the elastic cords. Lastly, the French Patent No. 74-04066 shows a chest exercising device which has three loops welded to a U-shaped support clamp. Three elastic members are attached to a handle and permit the user to exercise by pulling a handle away from the support brace.\nThe above devices all have useful provisions, but none of them may be safely used on a desk top while still being easily adjustable."}
{"text": "Pulse shape optimization has been suggested and developed for wireless and optical communication systems. However, conventionally, pulses are optimized according to metrics derived from a pulse's characteristics and each pulse shape was optimized for a specific objective such as, for example, tolerance to jitter, reduction of peak-to-average power ratio, energy maximization in high frequencies, improved tolerance to impulse truncation, etc. In fiber optic networks, pulse shaping is performed through filters at a transmitter and a corresponding receiver. Conventionally, optimization of the pulse shaping is done either theoretically (based on calculations) or based on an offline optimization procedure to minimize Intersymbol Interference (ISI) at the sampling phase of a receiver. Disadvantageously, these approaches do not take into account system impairments, such as finite quantification, filtering effects, the interaction between submodules, etc. Accordingly, conventional optimization approaches do not result in an optimal filter shape for an optical modem under different distortion conditions. Furthermore, conventional optimization approaches do not consider the fact that DWDM channels experience both intra- and inter-channel distortions. Therefore, obtained optimized pulse shapes, with conventional optimization approaches, do not result in a global optimal solution. Also, conventional optimization approaches have not been implemented as part of a working, multi-channel, end-to-end optical communication system."}
{"text": "The invention relates to a plug connector, especially for SMD (xe2x80x9csurface-mounted devicexe2x80x9d) plugs, for connection of electronic devices.\nPlug connectors are known and commonly used in diverse configurations. They have two plug-connector elements, which can be mutually engaged, and they are provided with a plurality of electrically conductive contacts that establish electrically conductive connections. The SMD plug connectors are plugged onto a printed-circuit board and detachably fixed by a snap fastening. Snap-in projections in the form of barbs engage in corresponding snap-in recesses. The plug-connector element disposed on the printed-circuit board can be detachably connected with another plug-connector element of complementary geometry. This permits the printed-circuit board to couple to a multi-conductor cable that has the other plug-connector element at its end.\nIt is known that the influence of external electromagnetic fields on a plug connector can be reduced by providing a shield plate of electrically conductive material on both plug-connector elements. This configuration also prevents electrical crosstalk in a plug connector. The shield plate surrounds their electrical contacts. These shield plates additionally function automatically as bonding conductors for grounding the plug connector or the printed-circuit board.\nUnder these conditions the problem of permanent electrical contact of the shield plates arises. This is especially true for plug-connector elements mounted permanently on the printed-circuit boards. In addition, the assembly process is considerably complex when a plug connector is to be equipped with a shield plate. Furthermore, a satisfactory solution has not yet been achieved for the electrical contact of these shield plates with one another.\nA person skilled in the art is faced with the object of designing a plug connector with shield plates in such a way that electrically conductive contact of the two shield plates with one another is permanently assured and that the shield plates can be contacted reliably with the associated printed-circuit board.\nThe central idea of the present invention is that, during assembly of a plug-connector element, the shield plates are automatically fixed via snap fastenings, so that displacement of the shield plate relative to the plug-connector element is prevented. It is ensured that, when the two plug-connector elements are coupled with one another, the respective shield plates bear directly against one another over substantially their entire surface. In other words, the shield plates are predominantly flush with their body areas facing one another, and do not slip. For this purpose the shield plates are constructed in the form of substantially flat metal components. Furthermore, soldering tags are formed on the shield plates, and project out of or extend beyond the connector-element body, which is usually a plastic injection-molded part. These soldering tags allow the plug-connector element with the shield plate to be permanently contacted with a printed-circuit board using a dip-soldering technique. The soldering tags are additionally used to position and align the shield plates during assembly of the plug-connector element and during plugging of the plug-connector element onto a wiring board. In the process the soldering tags extending beyond the plug-connector element are introduced into corresponding recesses of a wiring board or similar component, whereupon the element is automatically aligned in correct position. In addition, an appropriate retaining force for fastening the element is applied with a subsequently made soldered connection.\nThe shield plate and the soldering tags of the corresponding plug-connector element are connected in the same way to the bonding conductor of the multi-conductor electrical connecting cable.\nPreferably the plug-connector elements are constructed in the form of male and female multipoint connectors in each of which a plurality of electrical contact blades and sockets are disposed. Multipoint connectors generally have a rectangular body. These multipoint connectors are provided with snap-in projections or recesses, into which corresponding snap-in projections or recesses of the correspondingly shaped shield plate can be inserted to form a snap fastening. For this purpose the shield plate is preferably designed so that it comprises a continuous enclosure around the electrical contacts shielding the connection completely. The shield plate is at least along one of the side walls, preferably a long side of the plug-connector element. It can also have multi-piece construction.\nTo couple the male multipoint connector with the female multipoint connector, the male multipoint connector is provided with a receiving cavity into which the female multipoint connector can be introduced. For this purpose the blades of the male multipoint connector project from the base of the receiving cavity into the receiving cavity, and are contacted by the sockets arrayed in the female multipoint connector. The shield plates on the multipoint connectors are arranged so that the shield plate of the male multipoint connector is directly disposed on the inside of the receiving cavity and bears substantially against the body of the male multipoint connector, or the wall of the injection-molded part surrounding the receiving cavity. Correspondingly, the shield plate of the female multipoint connector is disposed on its exterior peripheral face. In the coupled condition of the two multipoint connectors, the flat and plane regions of the plates, facing one another, bear substantially completely and directly against one another.\nTo ensure that the two multipoint connectors are correctly positioned when plugged one into the other, guide elements are formed on the multipoint connectors. The guide elements have the form of projections, which extend outward beyond the body of the multipoint connector in plug-in direction and can be inserted into corresponding recesses on the other multipoint connector. This prevents the risk of skewed plugging of the two multipoint connectors and of bending of the blades or sockets.\nIt is also proposed that corresponding projections for guiding and aligning the multipoint connectors be provided on the shield plates, so that they also have one-piece construction.\nTo simplify plugging the two multipoint connectors one into the other, the edges or end regions of the shield plates pointing in the plug-in direction, are chamfered or are inclined relative to the body of the shield plate. The edges of the male multipoint connector are directed outward relative to the body and those of the female multipoint connector are directed inward relative to the body. Viewed in cross section, the portions each have symmetrical structure that either flares in the manner of a funnel or tapers in the manner of a pyramid. This facilitates the introduction of the one multipoint connector into the other even if the alignment is not exact.\nSMD plug-connector elements or leads are formed on the shield plates. This ensures that the shield plates are contacted and the electrical connections are shielded in the coupled condition. During assembly, therefore, the SMD lead of a shield plate of a plug-connector element is electrically contacted with, the printed-circuit board or wiring board onto which the plug-connector element is plugged."}
{"text": "A Surface Acoustic Wave (SAW) is an acoustic wave that propagates along the surface of a substrate. This principle has been commonly employed in SAW resonators. A SAW resonator is constructed with interdigital transducers (IDT) that provide a conversion medium between electric signals and surface waves. The SAW resonator is often constructed with an input IDT that generates a SAW from an electric signal, and an output IDT that detects the SAW traveling along the surface of the substrate. The frequency of the SAW resonator is based on the physical layout of the IDT and the velocity of the acoustic wave. The velocity of the acoustic wave is a function of the substrate material.\nSAW resonators are inherently high Q components, and therefore, make excellent, compact size, band pass filters for wireless communication devices. Their primary drawback, however, is their inability to be tuned by means of an external voltage variable capacitor. As a result, multi-band wireless handsets using SAW resonators typically require a separate channel for each operating band. Each channel includes transmit and receive filters coupled to a duplexer. Various switches are also needed to switch the selected channel to the antenna. As the number of desired bands increase, the cost and complexity of realizing a multi-band wireless handset with SAW resonators becomes prohibitive. Accordingly, there is a need in the art for a frequency agile SAW resonator. In wireless applications, frequency agile SAW resonators will enable wireless handsets to operate at multiple bands using a single channel. This will reduce power consumption, extend battery life, and reduce the size and cost of the wireless handset."}
{"text": "Liquid crystal (LC) materials are rod-like or lath-like molecules which have different optical properties along their long and short axes. The molecules exhibit some long range order so that locally they tend to adopt similar orientations to their neighbours. The local orientation of the long axes of the molecules is referred to as the “director”. Liquid crystals may broadly be classified as thermotropic (formed or modified by heat) or lyotropic (formed or modified by the action of a solvent). Commercial liquid crystal display devices use thermotropic liquid crystals.\nThere are three main types of thermotropic LC materials: nematic, cholesteric (chiral nematic), and smectic. For a liquid crystal to be used in a display device, it must typically be made to align in a defined manner in the “off” state and in a different defined manner in the “on” state, so that the display has different optical properties in each state. Two principal alignments are homeotropic (where the director is substantially perpendicular to the plane of the cell walls) and planar (where the director is inclined substantially parallel to the plane of the cell walls). In practice, planar alignments may be tilted with respect to the plane of a cell wall, and this tilt can be useful in aiding switching.\nHybrid Aligned Nematic (HAN), Vertical Aligned Nematic (VAN), Twisted nematic (TN) and super-twisted nematic (STN) cells are widely used as display devices in consumer and other products. The cells comprise a pair of opposed, spaced-apart cell walls with nematic liquid crystal material between them. The walls have transparent electrode patterns that define pixels between them.\nIn TN and STN displays, the inner surface of each wall is treated to produce a planar unidirectional alignment of the nematic director, with the alignment directions being at 90° to each other. This arrangement causes the nematic director to describe a quarter helix within the TN cell, so that polarised light is guided through 90° when a pixel is in the “field off” state. In an STN cell, the nematic liquid crystal is doped with a chiral additive to produce a helix of shorter pitch which rotates the plane of polarisation in the “field off” state. The “field off” state may be either white or black, depending on whether the cell is viewed through crossed or parallel polarisers. Applying a voltage across a pixel causes the nematic director to align normal to the walls in a homeotropic orientation, so that the plane of polarised light is not rotated in the “field on” state.\nIn a HAN cell, one wall is treated to align a nematic LC in a homeotropic alignment and the other wall is treated to induce a planar alignment, typically with some tilt to facilitate switching. The LC has positive dielectric anisotropy, and application of an electric field causes the LC directors to align normal to the walls so that the cell switches from a birefringent “field off” state to a non-birefringent “field on” state.\nIn the VAN mode, a nematic LC of negative dielectric anisotropy is homeotropically aligned in the “field off” state, and becomes birefringent in the “field on” state. A dichroic dye may be used to enhance contrast.\nLiquid crystal (LC) planar alignment is typically effected by the unidirectional rubbing of a thin polyimide alignment layer on the interior of the LC cell, which gives rise to a unidirectional alignment with a small pretilt angle. It has been proposed to increase the pretilt angle for a rubbed surface by incorporating small projections in the rubbed alignment layer, in “Pretilt angle control of liquid-crystal alignment by using projections on substrate surfaces for dual-domain TN-LCD” T. Yamamoto et al, J. SID, Apr. 2, 1996.\nWhilst having a desirable effect on the optical characteristics of the device, the rubbing process is not ideal as this requires many process steps, and high tolerance control of the rubbing parameters is needed to give uniform display substrates. Moreover, rubbing may cause static and mechanical damage of active matrix elements which sit under the alignment layer. Rubbing also produces dust, which is detrimental to display manufacture.\nPhotoalignment techniques have recently been introduced whereby exposure of certain polymer coating to polarised UV light can induce planar alignment. This avoids some of the problems with rubbing, but the coatings are sensitive to LC materials, and typically produce only low pre-tilt angles.\nAn alternative is to use patterned oblique evaporation of silicon oxide (SiO) to form the alignment layer. This also effects a desired optical response; however the process is complicated by the addition of vacuum deposition and a lithography process. Moreover, control of process parameters for SiO evaporation is critical to give uniformity, which is typically difficult to achieve over large areas.\nA useful summary of methods of aligning liquid crystals is given in “Alignment of Nematic Liquid Crystals and Their Mixtures”, J. Cognard, Mol. Cryst. Liq. Cryst. 1-78 (1982) Supplement 1.\nThe use of surface microstructures to align LCs has been known for many years, for example as described in “The Alignment of Liquid Crystals by Grooved Surfaces” D. W. Berriman, Mol. Cryst. Liq. Cryst. 23 215-231 1973.\nIt has been proposed in GB 2 286 467 to provide a sinusoidal bigrating on at least one cell wall, by exposing a photopolymer to an interference pattern of light generated by a laser. The bigrating permits the LC molecules to lie in two different planar angular directions, for example 45° or 90° apart. An asymmetric bigrating structure can cause tilt in one or both angular directions. Other examples of alignment by gratings are described in WO 96/24880, WO 97/14990 WO 99/34251, and “The liquid crystal alignment properties of photolithographic gratings”, J. Cheng and G. D. Boyd, Appl. Phys. Lett. 35(6) Sep. 15, 1979. In “Mechanically Bistable Liquid-Crystal Display Structures”, R. N. Thurston et al, IEEE trans. on Electron Devices, Vol. ED-27 No 11, November 1980, LC planar alignment by a periodic array of square structures is theorised.\nLC homeotropic alignment is also a difficult process to control, typically using a chemical treatment of the surface, such as lecithin or a chrome complex. These chemical treatments may not be stable over time, and may not adhere very uniformly to the surface to be treated. Homeotropic alignment has been achieved by the use of special polyimide resins (Japan Synthetic Rubber Co). These polyimides need high temperature curing which may not be desirable for low glass transition plastic substrates. Inorganic oxide layers may induce homeotropic alignment if deposited at suitable angles. This requires vacuum processes which are subject to the problems discussed above in relation to planar alignment. Another possibility for producing homeotropic alignment is to use a low surface energy material such as PTFE. However, PTFE gives only weak control of alignment angle and may be difficult to process.\nIt has been proposed, in EP 1 139 153 and EP 1 139 154, to provide a surface alignment structure comprising an array of alignment features which are shaped and/or orientated to produce a desired alignment. Depending on the geometry and spacing of the features the LC may be induced to adopt a planar, tilted, or homeotropic alignment. Such alignment structures have been used to provide bistable alignment in the Post Aligned Bistable Nematic (“PABN”) mode, as disclosed in EP 1 139 151.\nIn U.S. Pat. No. 5,155,610 and U.S. Pat. No. 5,262,882 it has been proposed to provide a surface layer of an anisotropic gel or polymer network containing non-reactive LC material on a substrate. The gel can effect an inclined orientation of LC molecules. The surface layer contains liquid crystal molecules, the orientation of at least a part of the molecules being permanently fixed in the anisotropic gel. The angle of inclination of the molecules of the layer differs maximally from a minimum at the interface with the substrate to a maximum at the interface with the LC material when the layer is sufficiently thick. The layer is applied in a thickness to produce a desired angle of inclination at the interface with the LC material.\nIn addition to controlling the alignment of LC molecules at cell wall surfaces, it is desirable to be able to control, or “tune” the anchoring energy at each surface. By doing this, the cell performance may be optimised for different display modes, for example HAN or VAN modes, or the bistable mode described in “Fast bistable nematic display using monostable surface switching”, I. Dozov et al, Appl. Phys. Lett. 70(9), Mar. 3, 1997. Here, different SiO thicknesses were used to provide different anchoring strengths. As discussed above, SiO evaporation is commercially undesirable for various reasons.\nIt is known from U.S. Pat. No. 5,073,294 to make a LC polymer film or fibre comprising polymeric liquid crystals having aligned multiple oriented mesogens. The products are said to be useful as high resolution imaging films, packaging films, filtration membranes, laminated films, waveguide devices and optical fibres."}
{"text": "Keeping power consumption at a low level is generally important for battery-powered devices, and particularly so for battery-powered wireless communication devices).\nUS 2011/0207509 A1 discloses various embodiments for providing enhanced battery power conservation in mobile devices. The device may store user pattern data and/or user preference data. The device may include a power management module coupled to the processor. The power management module may monitor a location of the device, wireless signal strength and date/time. The power management module may also monitor and analyze user operation of the device to identify scenarios in which wireless communications can be suspended to reduce battery consumption.\nUS 2014/0082384 A1 discloses techniques for power management of a portable device. According to one embodiment, a user agent of an operating system executed within a portable device is configured to monitor daily battery usage of a battery of the portable device, to capture daily battery charging pattern of the battery of the portable device, and to infer user intent of utilizing the portable device at a given point in time based on a battery operating condition at the point in time in view of the daily battery usage and the daily battery charging pattern. Power management logic is configured to perform power management actions based on the user intent.\nHowever, these power management techniques may not be adequate in all situations. For example, one or more of these power management techniques may cause power management that is either too restrictive or too indulgent in various situations.\nTherefore, there is a need for alternative power management for wireless communication devices."}
{"text": "It is known to mount the trim to a motor vehicle by driving a screw into a hole in a plastic nut mounted in the vehicle structure.\nIt is also known to have a chamfer edge on the hole in the plastic nut in order to break the sharp edges and help align the screw as it enters the hole. However, the area in which the screw can hit initially on the plastic nut and be guided to the hole is severally limited.\nIt would be desirable to provide a plastic nut which allows for mounting the trim with a blindly driven screw, where the plastic nut can direct the screw to one of a series of holes, when the screw first engages any portion of a planar trench."}
{"text": "1. Field of the Invention\nThe field of this invention is the construction of mobile units for the transportation of perishable food stuffs. More specifically the field of this invention is the construction of truck bodies adapted for the carring of hanging meats on tracks or rails within the body of the truck. Such tracks being adapted for releasably holding such hanging meats in a secure position within the truck. The field of this invention further embraces the means of releasing said securing means when the particular truck body is opened for the removal of such meats.\n2. Description of the Prior Art\nMany different prior art devices have been devised which teach various elements of the instant invention. Some of these various elements are typified by the following: U.S. Pat No. 2,590,533 by HAMER teaches generally the addition of cross beams for supporting the rails of overhead meat hanging structures thereby allowing a space between the roof and said cross beams, although it is not necessarily obvious therefrom that such space may be utilized for the storing of an overhead door. U.S. Pat. No. 3,572,815 by HACKNEY teaches the use of folding or overhead doors for truckbodies although not of the refrigerated variety and finally U.S. Pat. No. 2,827,118 by WENDT teaches insulated folding or overhead doors for refrigeration units. While these patents are typical of the prior art, there are several inherent problems left to be solved thereby. First, no means is taught showing how such hanging meats are prevented from sliding across the rails when some have been removed due to the acceleration of the carrying vehicle which still allow the hanging meats to be carried by the rails when removal is required. Further, no prior art devices allow the use of folding or overhead doors for refrigerated truck units.\nIt is an object of this invention to provide a truck body which allows overhead meat storage in combination with overhead doors.\nIt is a further object of this invention to provide a truck body wherein the hanging meats are securely held in relative position until removal is required."}
{"text": "There is a need for ever more accurate gyroscopes. In the prior art, Reference [1] below, which is incorporated herein by reference, describes gyroscope trimming for frequency tuning of the resonant modes of a gyroscope. In particular Reference [1] describes trimming to frequency tune a hemispherical resonating gyroscope (HRG), which is a large macro-scale cup sized structure with tines at the top lip of the wineglass. Tuning may be carried out either by machining the tines or chemically etching the tines. On the micro mechanical electrical system (MEMS) scale, tuning of gyroscopes has been explored both theoretically, as described in Reference [2] below, which is incorporated herein by reference, and experimentally, as described in References [3], [4] and [5] below, which are incorporated herein by reference.\nThere are several disadvantages with the approaches of the prior art. In almost all cases materials, such as metals, are initially added to a micro-scale structure, for example the structure described in Reference [3], so that the added material may be later ablated to tune the resonant frequencies. However, the added materials inherently lead to a low Q and frequency splits, which are the differences in resonant frequencies between one mode of resonance and another mode of resonance.\nFurther, where the region to be ablated is very close to the vibrating regions, trimming and therefore tuning may be difficult and very risky, as the structure may be ruined. Also the material properties of the structure may be degraded. In the case where tuning is carried out further away from the vibrating regions, surfaces are often not flat and not amenable to trimming.\nIn the prior art, several gyroscope resonating structures with spokes are described in references [6], [7], [8], and [9], below, which are incorporated herein by reference. The spoke structures of these prior art references generally have thicknesses that are similar to the thickness of the resonant structure, and/or are made of different materials, which leads to the disadvantages referred to above."}
{"text": "This invention relates to programs for digital computers, and more particularly to code translation for conversion of instruction code which was written for one computer architecture to code for a more advanced architecture.\nComputer architecture is the definition of the basic structure of a computer from the standpoint of what exactly can be performed by code written for this computer. Ordinarily, architecture is defined by such facts as the number of registers in the CPU, their size and content, the logic operations performed by the ALU, shifter, and the like, the addressing modes available, data types supported, memory management functions, etc. Usually, the architectural definition is expressed as an instruction set, and related elaboration.\nAs the technology used in constructing computers evolves, so does computer architecture. Semiconductor technology has served to make all structural features of a computer faster, less costly, smaller, lower in power dissipation, and more reliable. In view of such changes in the economics and performance of the computer hardware, it is necessary to make corresponding changes in architecture to take full advantage of existing hardware technology. For example, the CPU data paths have evolved from 16-bit, to 32-bit, to 64-bit. And, as memory has become cheaper, the addressing range has been greatly extended. A major departure in computer architecture, however, has been the retreat from adding more complex and powerful instructions, and instead architectures with reduced instruction sets have been shown to provide performance advantages.\nComplex instruction set or CISC processors are characterized by having a large number of instructions in their instruction set, often including memory-to-memory instructions with complex memory accessing modes. The instructions are usually of variable length, with simple instructions being only perhaps one byte in length, but the length ranging up to dozens of bytes. The VAX.TM. instruction set is a primary example of CISC and employs instructions having one to two byte /pcodes plus from zero to six operand specifiers, where each operand specifier is from one byte to many bytes in length. The size of the operand specifier depends upon the addressing mode, size of displacement (byte, word or longword), etc. The first byte of the operand specifier describes the addressing mode for that operand, while the opcode defines the number of operands: one, two or three. When the opcode itself is decoded, however, the total length of the instruction is not yet known to the processor because the operand specifiers have not yet been decoded. Another characteristic of processors of the VAX type is the use of byte or byte string memory references, in addition to quadword or longword references; that is, a memory reference may be of a length variable from one byte to multiple words, including unaligned byte references.\nReduced instruction set or RISC processors are characterized by a smaller number of instructions which are simple to decode, and by requiring that all arithmetic/logic operations be performed register-to-register. Another feature is that of allowing no complex memory accesses; all memory accesses are register load/store operations, and there are a small number of relatively simple addressing modes, i.e., only a few ways of specifying operand addresses. Instructions are of only one length, and memory accesses are of a standard data width, usually aligned. Instruction execution is of the direct hardwired type, as distinct from microcoding. There is a fixed instruction cycle time, and the instructions are defined to be relatively simple so that they all execute in one short cycle (on average, since pipelining will spread the actual execution over several cycles).\nOne advantage of CISC processors is in writing source code. The variety of powerful CISC instructions, memory accessing modes and data types should result in more work being done for each line of code (actually, compilers do not produce code taking full advantage of this). However, whatever gain in compactness of source code for a CISC processor is accomplished at the expense of execution time. Particularly as pipelining of instruction execution has become necessary to achieve performance levels demanded of systems presently, the data or state dependencies of successive instructions, and the vast differences in memory access time vs. machine cycle time, produce excessive stalls and exceptions, slowing execution. The advantage of a RISC processor is the speed of execution of code, but the disadvantage is that less is accomplished by each line of code, and the code to accomplish a given task is much more lengthy. One line of VAX code can accomplish the same as many lines of RISC code.\nWhen CPUs were much faster than memory, it was advantageous to do more work per instruction, because otherwise the CPU would always be waiting for the memory to deliver instructions--this factor lead to more complex instructions that encapsulated what would be otherwise implemented as subroutines. When CPU and memory speed became more balanced, a simple approach such as that of the RISC concepts became more feasible, assuming the memory system is able to deliver one instruction and some data in each cycle. Hierarchical memory techniques, as well as faster access cycles, provide these faster memory speeds. Another factor that has influenced the CISC vs. RISC choise is the change in relative cost of off-chip vs. on-chip interconnection resulting from VLSI construction of CPUs. Construction on chips instead of boards changes the economics--first it pays to make the architecture simple enough to be on one chip, then more on-chip memory is possible (and needed) to avoid going off-chip for memory references. A further factor in the comparison is that adding more complex instructions and addressing modes as in a CISC solution complicates (thus slows down) stages of the instruction execution process. The complex function might make the function execute faster than an equivalent sequence of simple instructions, but it can lengthen the instruction cycle time, making all instructions execute slower; thus an added function must increase the overall performance enough to compensate for the decrease in the instruction execution rate.\nThe performance advantages of RISC processors, taking into account these and other factors, is considered to outweigh the shortcomings, and, were it not for the existing software base, most new processors would probably be designed using RISC features. In order for software base, including operating systems and applications programs, to build up to a high level so that potential and existing users will have the advantages of making use of the product of the best available programming talent, a computer architecture must exhibit a substantial market share for a long period of time. If a new architecture was adopted every time the technology advances allowed it, the software base would never reach a viable level. This issue is partly alleviated by writing code in high level languages; a program written in C should be able to be compiled to run on a VAX/VMS operating system, or a UNIX operating system, or on MS/DOS, and used on various architectures supported by these operating systems. For performance reasons, however, a significant amount of code is written in assembly language, particularly operating systems, and critical parts of applications programs. Assembly language programs are architecture-dependent.\nBusiness enterprises (computer users, as well as hardware and software producers) have invested many years of operating background, including operator training as well as the cost of the code itself, in operating systems, applications programs and data structures using the CISC-type processors which were the most widely used in the past ten or fifteen years. The expense and disruption of /perations to rewrite the code and data structures by hand to accommodate a new processor architecture may not be justified, even though the performance advantages ultimately expected to be achieved would be substantial.\nCode translators are thus needed to ease the task of converting code written for one computer architecture to that executable on a more advanced architecture. The purpose of a code translator is to take in, as input, computer code written for execution on one type of architecture (e.g., VAX), and/or one operating system (e.g., VMS), and to produce as an output either executable code (object code) or assembly code for the advanced architecture. This is preferably to be done, of course, with a minimum of operator involvement. A particular task of a code translator is to detect latent error-producing features of the code, i.e., features that were acceptable in the prior use of the code as it executed on the previous operating system or architecture, but which may produce errors in the new environment."}
{"text": "This invention is related to a heater-humidifier-dehumidifier unit, and more particularly to a heater-humidifier-dehumidifier unit which is mounted in respective compartments of a housing.\nA lot of air conditioning apparatuses have been developed for providing proper living and working conditions in a defined area in proportion to the size of the particular apparatus. For example, a heater, a cooler, a humidifier and dehumidifier constitute known air conditioning apparatuses. These air conditioning apparatuses are separately manufactured, and each apparatus is chosen to provide a desirable condition in an area to be air-conditioned.\nTherefore, air conditioning apparatuses must be bought individually according to the need desired, that is, when warming of the area is desired, a heater is required. When removing moisture from humid air, a dehumidifier is required. When providing proper humidity in dry air, a humidifier is required. Thus, this causes an increase of installment costs, and an increase of dead space due to the occupation of a space by an unused air conditioning apparatus.\nFurther, the comfort of the space to be air-conditioned deteriorates in proportion to the operation of the air conditioning apparatus. That is, when in the warming operation, air becomes less humid or drier and when in a cooling operation, the air becomes more humid. Hence, a combination unit is required to compensate for the above problems."}
{"text": "The permanent alteration of the colour of keratinous fibres, in particular human hair, by the application of hair dyes is well known. In order to provide the consumer with the hair colour and the intensity of colour desired, a very complex chemical process is utilized. Permanent hair dyeing formulations typically comprise oxidative hair dye precursors, which can diffuse into the hair through the cuticle and into the cortex where they can then react with each other and suitable oxidising agents to form the end dye molecules. Due to the larger size of these resultant molecules they are unable to readily diffuse out of the hair during subsequent washing with water and/or detergents; hence delivering a consumer-desired permanency of colour. This reaction typically takes place in an aggressive environment at approximately pH 10 in the presence of an alkalizing agent and in the presence of an oxidizing agent. Moreover, the consumer repeats this process regularly in order to maintain the desired hair colour and shade and the intensity of colour and to ensure continual, even coverage of the hair including coverage of new hair growth.\nThe manufacturer of such products is also required to work within a large number of constraints. Since these products are being placed in direct contact with the consumers' skin, the potential exists for accidental contact with the eye or for ingestion (for example), which can occur during the dyeing process. Therefore, the formulation must meet rigorous safety requirements and not induce any allergic reactions. In addition to meeting these requirements, the products must also be optically and olfactory pleasing to the consumer. In particular, the products also need to meet certain physical parameters in order to ensure that the product can be easily applied to the hair by the consumer to provide the desired effect, without unintentional staining of the consumers' clothes, skin particularly along the hair line or other objects.\nThe manufacturer is also required to provide the hair colouring consumer a large range of different resulting colours. Some consumers may just wish to enhance the natural colour of the hair, whilst others may wish to cover grey or completely alter the hair colour to a different natural appearing hair colour or a ‘synthetic’ appearing hair colour. Consequently, the manufacturer may provide over twenty different formulations, of varying colours and shades, to address the range of consumer specific needs. These formulations have to be individually formulated and are typically complex formulae containing a mixture of different dye compounds. As a result the manufacture of such product ranges can be costly and complex.\nHowever, despite the fact that commercial hair dyeing products have been available for many years, the products still exhibit a number of consumer-related deficiencies.\nTypically permanent hair dye products will contain an alkali, typically a source of ammonia. This serves the purpose of swelling the hair allowing the entry of the dye precursor molecules into the hair and also improves the lightening effect of the oxidising agent, which is typically hydrogen peroxide. However, ammonia is also volatile and its associated odour is extremely unpleasant to the consumers' of such products, particularly as these hair dye products are used in close proximity to the nasal region. Hence, it would be highly desirable to provide an oxidative hair colouring and/or bleaching composition, which delivers the consumer required lightening level and colour, but which has reduced or eliminated the detectable ammonia odour.\nIn fact another deficiency area in current hair colouring products is the provision of hair colouring products which deliver the required hair lightening effect. Delivering the required level of lightening is particularly important in order to provide the full range of colour shades demanded by the consumer, especially for blonde shades and grey coverage. Such products pose particular difficulties to the manufacturer, as they usually require the use of high levels of oxidising agent and ammonia in order to deliver the required lightening effect. However, in additional to the problems associated with the presence of high levels of ammonia in these products, as discussed herein above, the presence of these high levels of ammonia and/or oxidizing agent also affect the condition of the hair and may in some cases induce mild skin irritation on the scalp. In particular, the hydrophilicity of the hair surface is increased during the colouring process, which alters the sensory perception of the hair and its overall manageability during, immediately after colouring and during the subsequent wash and styling cycles until the next colourant application. Hence, it would also be highly desirable to provide an oxidative hair colouring and/or bleaching composition which delivers the required lightening and/or colour without unnecessary hair damage.\nMoreover, in order to provide a product which the consumer can easily apply to the hair without dripping onto the skin, clothes or bathroom or salon surfaces, hair colourant products must be designed such that the applied composition has a certain required viscosity. This is either achieved by providing the dye composition and the oxidizing composition as so called thin-thin type liquid formulations which are thickened upon mixing, or, where at least one of the components, either the dye composition or the oxidizing composition, preferably the dye composition, is provided as a thickened formulation which thickens the total composition upon mixing. These thickened compositions can be achieved by the use of a gel network systems which provide the desired thickness to either the dye composition or the oxidizing composition or, preferably both compositions. Furthermore, such gel networks are highly desirable as they also provide additional benefits of a cream like texture, conditioner like feel and appearance, smooth rinse and improved hair feel. Such thickened gel network systems are described in WO2007/102119, EP1878469 and EP1878469. These compositions have a viscosity range typical for retail applications which is in the range of 1 Pas to 8 Pas.\nIn the professional hair salon market, the hair colour professionals typically do not use predetermined combinations of dye and oxidizing compositions or kits as in the retail market. Instead, the hair salon professionals are able to use any combination of commercially available dye and oxidizing compositions. For the professional market, the available oxidizing compositions are typically provided with a variety of viscosities. Consequently, the dye compositions that are combined with the oxidizing compositions need to be formulated to ensure that they are able to thicken the oxidizing compositions to the required viscosity independent of the viscosity of the starting composition of the oxidizing agent.\nIn addition, in the professional salon the mixed dye and oxidative compositions are typically much thicker than the mixed compositions provided for the retail market and are typically in the range of 9 Pas to 16 Pas. This is to enable the hair salon professionals to use a variety of techniques to apply the composition in order to provide their clients with the desired results. Typically this is delivered using the brush and bowl technique. Compositions which have a viscosity typical to retail applications are considered not acceptable by hair colour professionals.\nHowever, certain professional thickened compositions comprising such gel network systems have been found to exhibit a phenomenon whereby after application to the hair, the composition tends to crack away from the root along the length of the hair strand away from the root. As a result the roots do not necessarily receive the required amount of composition in order to deliver the desired result. Particularly for regular hair colourant users, virgin hair growth coverage at the roots is essential for a satisfactory outcome and this problem needs to be addressed, without the necessity for the hair salon professional to continuously check and reapply the composition to the roots to ensure complete coverage. This problem of lack of effective root adhesion is further exacerbated with compositions comprising hydrogen peroxide and carbonate oxidizing systems.\nHence it would be desirable to provide the consumer, particularly the hair salon professional with a hair colorant product, which in addition to delivering the required lightening, colour deposition, has the required rheology and viscosity such that it can be utilised by hair salon professionals and which does not exhibit any cracking away from the roots during or after application of the composition on the hair.\nIt has now been surprisingly found that oxidative hair colouring compositions comprising at least one oxidizing agent, a specified gel network thickener system comprising a specified tertiary system, and a cationic polymer and a mica or titanium oxide as defined herein having the required viscosity and ionic strength so it can be used to achieve variety of professional hair colour applications while exhibiting excellent hair root adhesion."}
{"text": "Light microscopes, transmission electron microscopes, scanning electron microscopes and other instruments are extensively used to understand the ultrastructure of a wide variety of synthetic and biological materials in numerous areas of science and technology. For example, light microscopy samples are used for research to identify the development of different organs in animals and plants. In addition, one major use of LM samples is in the histophathologic examination of biopsy samples of tissues suspected of disease. TEM is used to study both biological samples as well as non-biological samples and can even provide atomic scale resolution. TEM is routinely used to investigate metal grain structures, the micro and nano-structures of polymers, semiconductor devices such as computer chips, and to visualize the organelle and molecular organization of cells. Such images are capable of resolution down to approximately 0.1 nm, although this is usually not quite possible with biological materials. TEM and LM specimens are commonly sliced or otherwise prepared into thin cross-sections to enable electrons or photons, respectively, to traverse through the specimen to create an image. SEM is similar to TEMs in that it uses electrons to create an image of the target/sample. However, the resolution of the SEM is typically not as fine as that of the TEM, yet high resolution versions are capable of molecular level resolution of approximately 5 nm. Due to the SEM's ability to image the surface aspect of bulk materials, specimen preparation does not usually entail slicing the specimen into cross-sections.\nStudy objects for microscopy are prepared in multiple ways depending upon the type of material to be examined, and the type of microscopy to be used. Biological materials require special handling to preserve the structure of the material when it will be examined in the electron microscope, and secondarily to enhance or enable imaging.\nBoth SEM and TEM instruments perform their imaging in a vacuum (the absence of air or other gases). Since biological materials are generally 50 to 95% water, if these were placed directly within the vacuum the water would evaporate and the specimen would collapse. Consequently both SEM and TEM samples have the water removed after the structure is strengthened with chemicals such as glutaraldehyde, formaldehyde, and osmium tetroxide.\nTEM samples must be very thin (typically about 40-100 nm) in order for the electrons used for imaging to be “transmitted” or passed through the sample. To cut specimens into such thin sections the water in the specimen is replaced with plastic that is hardened in place. This plastic supports the sample as it is sliced very thin using a device called an ultra-microtome.\nLM specimens, especially those of biological origin, are generally also sectioned in order to provide cross-sections for viewing, and to allow photons (light) to be transmitted through the specimen. As with TEM, LM sections are also embedded to support the specimen during sectioning, however, different generally softer plastics or wax are used as the embedding materials. Water that is frozen with additional materials to enable the ice to be softer, provide better support of the tissue, and reduce ice crystal damage during freezing can also be used.\nBecause of the delicate nature of microscopy samples and the resolution and power of LM, SEM and TEM analysis, sample preparation requires highly skilled and delicate manipulation. For example, the preparation of LM, SEM and TEM samples may require from between 20 to 40 fluid exchanges. Such samples may be prepared in large quantity and their quality, analysis and identification may have significant downstream impact, such as for biopsy samples, for example. Thus, methods and devices for their preparation have been devised that are either highly cumbersome and/or inadequate for their needs.\nFor example, U.S. Pat. No. 5,080,869 to McCormick discloses an apparatus and method for preparing samples for histological examination. The device comprises a cassette suitable for holding a sample with perforation in the wall and floor to drain fluid. A stack of such cassettes can be put in a container for the process of fixation. As illustrated by McCormick, inside the container and the cassette are large amounts of dead space resulting in the need for large volumes of fixation media, rinses, and other solutions in order to adequately process the sample. Since many of the chemicals used are noxious and some are expensive to purchase or dispose of, large dead-space volumes are not desirable. In addition, such cassettes are not intended for the preparation and handling of the circa 1 mm specimens typical in biological TEM and many clinical histopathology or biopsy specimens and hence such small specimens can readily become lost or misplaced.\nU.S. Pat. No. 5,543,114 to Dudek discloses a unitary biological specimen processing apparatus. The apparatus comprises a container having a lid with apertures in it for straining fluid from the container. As disclosed by Dudek, a sample is placed in the container and fixation fluids entered in one end and emptied out the other end. As can be seen in the Dudek figures, the container has a large volume of dead space in which the sample can be lost or damaged. Further, the sample can not be visually inspected nor can the sample be sectioned or stored in the container.\nSimilarly, U.S. Pat. No. 7,179,424 to Williamson et al. discloses a cassette for handling and holding tissue samples during processing. As with other, similar devices, the cassette disclosed by Williamson includes a large volume of dead space, and provides little ability to visually inspect the samples within. Further, the cassettes taught by Williamson are not amendable to simultaneous use with other cassettes thus limiting their ability for high-throughput use. In addition, as taught by Williamson et al., the cassettes are complicated devices having relatively high costs and not allowing for sectioning of the sample held within.\nFurther, most other imaging and analytical methods used to analyze specimens on such small scales could benefit from methods and/or devices that facilitate and standardize the process of preparing such samples. Such applications require the handling, processing, and identification of small specimens for analyses such as, for example, secondary ion mass spectroscopy (SIMS), electron spectroscopy for chemical analysis (ESCA) which is also known as x-ray photoelectron spectroscopy (XPS) and atom probe tomography. Matrix-assisted laser desorption ionization (MALDI), and many other analytical and imaging instrument preparation procedures, are also within the scope of this invention.\nThus, the need exists for a low-cost device and method that allows for the parallel uses of fixation of multiple samples, providing for a high-throughput fixation system that further allows for the tracking and identification of samples held both on a short-term scale and for long-term storage.\nTherefore, there is a need for devices and methods to be used in the preparation of TEM grids and specimens for TEM, SEM, LM, SIMS, ESCA, XPS and MALDI that provide for easier preparation of the grids and specimens for analysis and that may allow for more efficient storage and handling."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for the biological preparation of 25-hydroxyvitamin D compounds and 1.alpha.,25-dihydroxyvitamin D compounds from vitamin D compounds.\n2. Prior Art\nVitamin D compounds are known to be important vitamins which participate with absorption of calcium in the body and acceleration of calcium metabolism of the bone, and to be activated by hydroxylating at the 25-position in the liver and at the 1.alpha.-position in the kidney [H. F. DeLuca and H. K. Johns, in Ann. Rev. Biochem., vol. 52, page 411 (1983)]. However, patients with liver function disorders are deficient in production content of 25-hydroxyvitamin D.sub.3 in the body. Accordingly, administration of 25-hydroxyvitamin D.sub.3 to these patients is effective from the viewpoint of a supply of the vitamin deficiency in the body. Furthermore, 1.alpha.,25-dihydroxyvitamin D.sub.3 has especially strong activity, and shows remarkable effect on patients with renal failure from the same viewpoint as described above. It is known that 1.alpha.,25-dihydroxyvitamin D.sub.3 can be obtained from 25-hydroxyvitamin D.sub.3 as a material by biological and enzymatic chemical conversions (Japanese Patent Kokai 2-469 and ibid. 2-231089). Accordingly, it is industrially advantageous to prepare 25-hydroxyvitamin D.sub.3 and 1.alpha.,25-dihydroxyvitamin D.sub.3 inexpensively.\nDirect introduction of hydroxyl groups into vitamin D.sub.3 compounds at the 1.alpha.- and/or 25-positions by organic synthesis is not known. Such an introduction by using animal organs is known [Biochem. Biophys. Res. Commun., vol. 36, page 251 (1969)], but, this method is not industrially practical. Hydroxylation of vitamin D.sub.3 at the 1.alpha.- and/or 25-positions by enzymatic chemistry using microorganisms was recently disclosed (Japanese Patent Kokai 2-469 and ibid. 2-231089).\nAccording to the preparations of 25-hydroxyvitamin D compounds and 1.alpha.,25-dihydroxyvitamin D compounds by enzymatic chemistry using the microorganisms disclosed in Japanese Patent Kokai 2-469 and ibid. 2-231089, the ratio of conversion from the substrate vitamin D.sub.3 and 25-hydroxyvitamin D.sub.3 into 25-hydroxyvitamin D.sub.3 and 1.alpha.,25-dihydroxyvitamin D.sub.3, respectively, is not satisfactory. Furthermore, these methods have problems to solve for the industrial preparation, for example, since the concentration of the substrate used is low, these methods need large scale equipment for industrial production.\nThus, an object of the present invention is to provide an improved method for the biological preparation of 25-hydroxyvitamin D compounds and 1.alpha.,25-dihydroxyvitamin D compounds from vitamin D compounds and 25-hydroxyvitamin D compounds, respectively, using the prior art microorganisms or enzymes with remarkably increased conversion ratio.\nAs a result of earnest research in order to increase the ratio of conversion from vitamin D compounds and 25-hydroxyvitamin D compounds into 25-hydroxyvitamin D compounds and 1.alpha.,25-dihydroxyvitamin D compounds, respectively, by using microorganisms, the present inventors have found that the conversion ratio is remarkably increased by adding cyclodextrin compounds to a reaction mixture."}
{"text": "In cases where individuals have experienced sensorineural deafness, the restoration of hearing sensations to such individuals has been achieved through the use of hearing aids and cochlear implants. Cochlear implants in particular have been in clinical use for many years. A wide variety of different speech processing strategies have been employed in order to process a sound signal into a basis for electrical stimulation via implanted electrode arrays. Some systems have focused upon extracting particular acoustic components of the detected sound signal, which are important to the user's understanding of speech, for example the amplitudes and frequencies of formants, and using these as a basis for generating stimuli. Other approaches have also attempted to utilise the generally tonotopic arrangement of the cochlea, so that each electrode corresponds generally to a particular frequency band.\nOne such approach, commercially used in speech processors sold by Cochlear Limited, is known as SPEAK. In the SPEAK system, the incoming sound signal is processed to provide an indication of the amplitude of the ambient sound signal in each of a predetermined set of frequency channels, and the channels with the largest amplitudes are selected as the basis for stimulation. In other approaches, the outputs of all channels are used to specify the stimulation patterns, rather than just the channels having the highest short-term amplitudes. The channels are defined by the partially overlapping frequency responses of a bank of band-pass filters. The filters may be implemented using a variety of analog or digital techniques, including the Fast Fourier Transform (FFT). The electrodes corresponding to those channels, determined by a clinical mapping procedure, are selected for activation in each stimulation period and are allocated to the channels according to the tonotopic organization of the cochlear. The rate of stimulation is preferably as high as possible subject to limitations imposed by the processing and power capacity of the external processor and implanted receiver/stimulator unit.\nThe range of electrical stimulus levels is usually determined by psychophysical measurement of threshold and comfortably loud levels on individual electrodes, using fixed-current pulse trains at the same rate as the stimulus cycle rate of the speech processor output. This may be described as per electrode loudness mapping. The problem with this method of loudness-mapping is that is does not take into consideration the effects of loudness summation when multiple electrodes are activated in quick succession, as they generally are in the output of speech processors.\nAlthough most processing strategies activate a nominal fixed number of electrodes per stimulus cycle, it is important to realise that the actual number of electrodes stimulated in individual cycles is a variable subset of this number, depending on the level and bandwidth of the acoustic stimulus at each point in time. To illustrate this point, a low-level acoustic pure tone will lead to activation of a single electrode, and the electrical level on this electrode must be at least equal to the psychophysical threshold measured individually for that electrode to be audible. In contrast, a low level broad-band noise may activate (for example) eight electrodes in a stimulus cycle. If each of these eight electrodes are activated close to their individual psychophysical thresholds, as may occur with existing systems, then the resultant loudness will not be close to threshold loudness as intended, but will be closer to the maximum comfortable loudness.\nThis loudness summation leads to the situation that the output of the processor is too loud, even though the individual levels on each electrode do not exceed a comfortable loudness. Various practical methods have been employed to attempt to overcome this problem, including a global reduction of the upper level limit on each electrode, or the use of complex input signals to set the range of individual levels across electrodes. These methods, although alleviating the discomfort of implant users for loud sounds, do not address a second important issue, and that is the impact of loudness summation on speech perception.\nAmplitude envelope fluctuations of a speech signal provide vital cues for speech perception, especially for those people who are less able to make use of spectral cues in the signal (for example, those with few active electrodes or poor electrode discrimination ability). Therefore it is important that the changes in acoustic intensity from moment to moment in a speech signal are accurately conveyed as the appropriate perceptual loudness changes to the implantee. The present loudness coding methods, whereby the acoustic output of a filter is mapped to a fixed range of electrical levels (however determined) on its corresponding electrode, lead inevitably to a perceptual distortion of the amplitude envelope shape because these methods do not take into account the variations from moment to moment of important aspects such as the number of electrodes activated in each stimulus cycle, and the relative loudness contributions from these other electrodes. In summary, the relative loudness of electrically stimulated hearing using present approaches does not accurately convey the relative loudness that a normally-hearing person would hear for the same acoustic input. As well as distorting the perception of the amplitude envelope of the acoustic signal, this effect will lead to narrow-band signals being masked by lower-level broad-band noise, thus disrupting the ability of implantees to understand speech in background noise.\nWhilst the SPEAK approach has proven successful clinically, it is an object of the present invention to improve sound processing strategies so as to enhance intelligibility of speech and other sounds, for users of cochlear implants. It is a further object of the present invention to improve the perception of loudness provided to users of cochlear implants."}
{"text": "1. Technical Field\nThe present invention relates to an image forming apparatus, in which a toner image formed on a transfer member, based on image data, is transferred and fused on a recording medium to form an image.\n2. Description of the Background Art\nIn electrophotographic image forming apparatuses, toner images are formed on a transfer belt using image data, and then transferred and fused on a sheet to form images. Such image forming apparatuses can form images using colored toner such as yellow, magenta, cyan, and black toner, and also clear toner. For example, an image forming apparatus can form images having a watermark on the top layer of the images by superimposing a toner image of clear toner over the toner images of colored toner, and transferring and fusing each of toner images on a sheet.\nCompared to color image formation using only the colored toner, when an image is formed using a combination of the colored toner and the clear toner, the total amount of toner to form the image becomes great. However, in image forming apparatuses, if the total amount of toner transferred on the sheet becomes great, heat amount used for plasticizing or melting toner for the normal transfer process may not be enough for fusing a toner image on the sheet, and resultantly a transfer failure may occur.\nIn light of such problem, JP-2009-63744-A discloses an image forming apparatus using an image forming method in which an upper limit (or control value) is set for the total combined amount of colored toner and clear toner, and when the total amount of toner exceeds the upper limit, the density or concentration of colored toner is adjusted. In such a method, when the total amount of toner exceeds the upper limit, the density or concentration of clear toner is fixed to a given value, and the density or concentration of colored toner is decreased to limit the total amount of toner at the upper limit. With such a configuration, an image can be formed without affecting the gloss provided by the clear toner.\nHowever, such toner-amount reduction of the colored toner amount has the undesirable effect of decreasing image density of the resultant output image."}
{"text": "1. Field of the Invention\nThis present invention relates generally to the field of integrated circuit design and, more specifically, to an input buffer circuit that has operational characteristics that are selectable depending on the intended application.\n2. Description of the Related Art\nThis section is intended to introduce the reader to various aspects of art that may be related to various aspects of the present invention, which are described and/or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present invention. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nComputer systems and other electronic devices typically include a variety of electrically interconnected integrated circuit (IC) packages which perform a variety of functions, including memory and processing functions. Many integrated circuit devices have input buffers, which receive data from outside of the integrated circuit. The input buffer is the first stop for data inside the integrated circuit before the data is stored or further processed. Typically, one input buffer is used for each data input of an integrated circuit.\nThere are two commonly used types of input buffers: (1) differential input buffers and (2) Complementary Metal Oxide Semiconductor (CMOS) input buffers. Each of these input buffer types has advantages and disadvantages compared to the other.\nA differential input buffer is typically used to monitor the transitions of an input signal relative to a fixed reference voltage or relative to a compliment differential signal. The fixed reference voltage may be referred to as VREF. The output of a differential input buffer may transition from a logical high (xe2x80x9c1xe2x80x9d) to a logical low (xe2x80x9c0xe2x80x9d) or vice versa as the input crosses the threshold established by the reference voltage VREF. Differential input buffers typically offer higher performance (faster switching time) than CMOS buffers, but differential input buffers also tend to consume more power than CMOS input buffers. The high power consumption of differential input buffers is in part because differential input buffers have source voltages that are always present, even when the buffer is not in normal operation. These source voltages are partially dissipated by the internal circuitry of the differential buffer, which consumes power. Also, the reference voltage VREF is typically always connected to a differential input buffer, resulting in further power consumption.\nCMOS input buffers are simpler than differential input buffers. CMOS input buffers do not require the connection of a reference voltage VREF. Also, the internal circuitry of CMOS buffers does not consume power when the buffer is not in operation. As a result, CMOS buffers do not generally consume as much power as differential input buffers. A disadvantage of CMOS input buffers is that they do not typically offer performance levels in terms of switching speed associated with differential input buffers.\nPower consumption and performance (operational speed) are two important factors in the design of integrated circuits. Accordingly, those factors are important in the design of input buffers, as well. The relationship between power consumption and performance is a classic design tradeoff. This means that an input buffer that is optimized for performance (high speed) typically consumes much more power than an input buffer designed to save power. Correspondingly, an input buffer that is designed to conserve power typically has slower performance characteristics compared to an input buffer that is designed for high performance.\nAs new designs for computer systems and other electronic devices continue to evolve and proliferate, designers of electronic devices are constantly faced with deciding what performance characteristics they need to incorporate in devices having different operational characteristics. For example, a small portable electronic device such as a cell phone or a personal digital assistant (PDA) may have reduced power consumption as a primary design goal. On the other hand, a high performance desktop computer system that is plugged into a source of AC power most of the time may have higher performance as a more important design goal than low power consumption.\nMany computers and electronic devices employ volatile random access memory (RAM) to store information during normal operation. One examples of a type of RAM that is commonly used in a wide range of electronic devices is dynamic random access memory (DRAM), which must be continually updated or xe2x80x9crefreshedxe2x80x9d to preserve its contents. DRAM is typically used for the main memory of most computer systems. There are many different implementations and types of DRAM, including synchronous DRAM (SDRAM) and double data rate (DDR) SDRAM.\nThe power consumption and performance tradeoff discussed above is applicable to the choice of input buffers for the DRAM memory in a given electronic device design. This is true because DRAM can have a major impact on the overall power budget of an electronic device. Additionally, the performance of DRAM may have significant impact on the overall performance of an electronic device. In the past, designers of electronic devices may have had to choose a DRAM with CMOS input buffers for certain applications (e.g. low power consumption) and a different DRAM configuration or implementation with differential input buffers for other applications (e.g. higher performance). Specifically, a designer of a cell phone or other small electronic device may have chosen a DRAM with CMOS input buffers to meet a design requirement of low power consumption while a designer of a desktop computer system may have had to choose a different type of DRAM with differential input buffers to meet a design requirement of high system performance. A DRAM having operational characteristics that may be selected depending on the specific application is desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to point-to-point or point-to-multipoint wireless communications systems. More specifically, the present invention relates to systems and methods for sharing various forms of media between members of a communicating group of wireless telecommunication devices.\n2. Description of the Related Art\nIn wireless telecommunication devices, such as cellular phones, PDAs, mini-laptops, and advanced pagers, the devices typically communicate over long distances by bridging telephone calls through existing cellular telephone networks and passing data packets across the network. These wireless devices often have limited to significant data processing and computing capability, and can accordingly send and receive software programs, in addition to voice, across the telephone network.\nThere exists a wireless telecommunication service that provides a quick one-to-one or one-to-many communication that is generically referred to as “Push-To-Talk” (PTT) capability. The specific PTT group of recipient devices for the communicating wireless device is commonly set up by the carrier. A PTT communication connection is typically initiated by a single button-push on the wireless device that activates a half-duplex link between the speaker and each member device of the group and once the button is released, the device can receive incoming PTT transmissions once the button is released. In some arrangements, the PTT speaker will have the “floor” where no other group member can speak while the speaker is speaking. Once the speaker releases the PTT button, any other individual member of the group can engage their PTT button and they will have the floor.\nIn existing PTT systems, member devices will only communicate voice data between the devices in an active group communication. It is difficult to include other data in the group communication because of the size of the data and the device resources necessary to handle non-voice data. The existing mobile devices therefore send non-voice data, such as graphics, multimedia and other data files, through specific data channels established between the device and a wireless communication network.\nAccordingly, it would be advantageous to provide a system and method for a wireless device member of a group to send media to others member of the PTT target group. The system and method should be able to accommodate group-directed media data of large size without compromising the integrity of other PTT communications or unduly utilizing device resources. It is thus to the provision of such a system and method to allow group-directed media among wireless devices in a PTT group that the present invention is primarily directed."}
{"text": "Manually operable, molded plastic dispensers, that are secured to containers for diverse products, such as hair care products, beauty aids, skin creams lotions, health care products, etc. are well known and have gained widespread acceptance. Such dispensers must be inexpensive to manufacture, and assemble, lend themselves to installation on containers by automated machinery operating at high throughput, and function satisfactorily over the life span of the product being dispensed. Also, such dispensers must be aesthetically pleasing, capable of being sealed to avoid leakage during shipment, and capable of being manipulated between an opened, and a closed, condition, without the accumulation of unsightly excess product on exposed surfaces. A representative dispenser that addresses some, if not all, of the concerns noted above, is disclosed in U.S. Pat. No. 5,472,120 granted Dec. 5, 1995 to Stebick et al.\nStebick et al discloses a cap 10 for a container having three distinct positions, namely, a closed position shown in FIG. 1A, an intermediate spray position shown in FIG. 1B, and a fully open, pour position shown in FIG. 1C. The cap provides a fluid tight seal when in the closed position, as indicated by sealing areas A, B, C, and D, shown in FIGS. 1A and 4. The cap provides a fluid tight seal from pouring, when in the spray position, for plug 38 on stem 32 of chimney 24, fits snugly within central pouring aperture 54.\nCap 10 in Stebick et al includes a shell 12 defining a fluid passage 42 to the container, and a tip 14 with pouring aperture 54 and spray 56 apertures axially movable relative to the shell. The extent of the axial movement is governed by the engagement of annular rings 64 on the inner surface of sidewall 52 of tip 14, as shown in FIG. 3C, with the axially spaced grooves or indents 28 in the outer wall of chimney 24, as shown in FIG. 2A. Tip 14 is pulled up to select a dispensing position, which causes plug 38 to engage, or disengage, with aperture 54, depending on the selected dispensing option."}
{"text": "The present invention relates generally to an electrodeposited copper foil with its surface prepared wherein a matte side, namely a surface at which copper deposition is completed, is mechanically polished at least once and the thus mechanically polished matte side is subjected to a selective chemical polishing so as to attain preparing thereof, and relates to a process for producing the same and a use of the electrodeposited copper foil with its surface prepared in, for example, a printed wiring board.\nIn recent years, both the size and weight of electronic equipment such as a notebook-sized personal computer are being reduced to increasing extents. Accordingly, IC wiring is also becoming finer.\nWith respect to the wiring pattern formed on a substrate used in such an electronic equipment, the lead width is now as small as ten-odd microns (xcexcm). In accordance therewith, the metal foil constituting the wiring pattern is becoming thinner. Specifically, while the designated thickness of metal foil for use in the formation of the conventional wiring pattern of about 100 xcexcm lead width has ranged from about 15 to 35 xcexcm in correspondence to the width of the wiring pattern, the thickness of metal foil employed in the formation of ten-odd micron (xcexcm) wiring pattern must be reduced in correspondence thereto.\nFor example, an aluminum foil or a copper foil is used as the metal foil for constituting the above wiring pattern. It is preferred to employ a copper foil, especially an electrodeposited copper foil, as the metal foil.\nThe electrodeposited copper foil employed for forming the above wiring pattern is produced by electrodepositing copper on a drum surface. With respect to the thus produced electrodeposited copper foil, the surface at which copper deposition is initiated, namely the surface at which formation of copper deposits brought into contact with the drum is initiated, is referred to as xe2x80x9cshiny sidexe2x80x9d, and the surface at which copper deposition is completed is referred to as xe2x80x9cmatte sidexe2x80x9d. The surface condition of the shiny side is substantially the same as that of the drum. That is, the 10-point average surface roughness (Rz) in ISO 4287 of the drum is about from 1.2 to 2.5 xcexcm, to which the 10-point average surface roughness (Rz) of the shiny side is nearly equal. On the other hand, with respect to the matte side, its surface roughness is greater than the surface roughness of the shiny side, and the 10-point average surface roughness of the matte side, although varied depending on the deposition condition of copper and the thickness thereof, is generally in the range of about 2.5 to 10 xcexcm. In the conventional electrodeposited copper foil of about 35 xcexcm nominal thickness, it has been rare that the surface roughness of the matte side poses a problem. However, in the electrodeposited copper foil of ten-odd micron (xcexcm) thickness, the surface roughness of the matte side is equivalent to tens of percents of the thickness of the whole electrodeposited copper foil, and the condition of the matte side exerts marked influence on the electrical properties of formed wiring pattern and board per se. It is known that, for example, mechanical polishing, chemical polishing and electrolytic polishing are available as the means for preparing the state of surface of the copper foil. The mechanical polishing is a method of smoothing the surface of the copper fail with the use of, for example, a buff. When use is made of a thin copper foil, the copper foil may be broken by mechanical stress exerted on the copper foil. Thus, the mechanical polishing is suitable for the preparing of the surface of relatively thick copper foils. On the other hand, no mechanical stress is exerted on the copper foil in the chemical polishing and electrolytic polishing, as different from the mechanical polishing, so that even thin copper foils would not be broken by the chemical polishing and electrolytic polishing. Thus, it has been believed that the chemical polishing and electrolytic polishing are suitable for the preparing of the surface of relatively thin copper foils.\nFor example, Japanese Patent Unexamined Publication No. Hei 5-160208 discloses a tape carrier having a lead pattern formed from an electrodeposited copper foil wherein the overall surface of matte side obtained by electrodeposition has been polished. This publication discloses the use, in the formation of a lead pattern of 60 to 80 xcexcm pitch, of an electrodeposited copper foil whose 1-2 xcexcm of matte side surface profile has been chemically polished to thereby attain preparing of the matte side. The thickness of the there employed electrodeposited copper foil after the preparing is in the range of 18 to 30 xcexcm. It is disclosed that a highly reliable carrier tape with desired lead strength can be provided by the use of the copper foil whose matte side overall surface has been chemically polished.\nHowever, the preparing of copper foil by chemical polishing as described in the above publication, although protrudent parts of the matte side are leached with relatively high selectivity to thereby effect preparing thereof, also invites leaching of the copper constituting the depressed parts of the matte side. Therefore, in this chemical polishing, the whole copper foil tends to become thin. Accordingly, when the thin electrodeposited copper foil employed in conformity with the recent trend toward fine pitch, for example, the electrodeposited copper foil having a thickness of 35 xcexcm (1 ounce), or 17.5 xcexcm (xc2xd ounce), or less is chemically polished to such an extent that a desired state of surface is attained, the whole electrodeposited copper foil is thinned to such an extent that the mechanical strength of wiring pattern or lead becomes poor. Further, this chemical surface polishing poses a problem such that it is difficult to control reaction conditions for chemical polishing so as to have the matte side uniformly treated. These problems of chemical polishing also occur in the electrolytic polishing involving leaching of copper.\nIn this connection, Japanese Patent Application Publication (Unexamined) No. Hei 3-296238 discloses a method of producing a TAB tape having a wiring pattern formed from a non-treated copper foil. The average surface roughness of the non-treated copper foil is described as falling within the range of 0.01 to 1 xcexcm.\nHowever, the non-treated copper foil whose average surface roughness (Rz) falls within the range of 0.01 to 1 xcexcm, disclosed in this publication, is a rolled copper foil. The surface roughness of this non-treated rolled copper foil is too low to ensure satisfactory peel strength (bonding strength). Accordingly, it is needed to preheat the copper foil or increase the diameter of the roller so as to form a covering thin film of cuprous oxide on the surface of the rolled copper foil. This poses a problem such that the process becomes laborious. Further, the use of this rolled copper foil renders it difficult to form a wiring pattern of extremely fine pitch such as one of from 30 xcexcm to less than 60 xcexcm pitch width.\nStill further, Japanese Patent Unexamined Publication No. Hei 9-195096 discloses an invention directed to an electrodeposited copper foil for printed wiring board wherein the surface roughness (Rz) of the matte side of electrodeposited copper foil prior to nodulating treatment is not greater than 1.5 xcexcm while the surface roughness (Rz) after nodulating treatment on the matte side is in the range of 1.5 to 2.0 xcexcm. This electrodeposited copper foil is described as being producible by a method comprising buffing the matte side of an electrodeposited copper foil so as to cause the surface roughness (Rz) prior to roughening treatment to become in the range of 1.5 xcexcm or less and subsequently effecting a roughening treatment on the matte side so as to cause the surface roughness (Rz) to become in the range of 1.5 to 2 xcexcm.\nHowever, the buffing of electrodeposited copper foil at a stretch as described in this publication may cause streaks on the buffed surface. These streaks result from polishing made deeper than predetermined. Some streaks have not posed any problem when use is made of the conventional thick electrodeposited copper foils. However, these streak portions are formed by excess polishing of copper, so that, when use is made of thin copper foils, the depth of streak portions is extremely large. Thus, these are likely to become the cause of defective occurrence, for example, high possibility of open circuit at such portions in a wiring pattern or the like. Furthermore, in the execution of such buffing, stress is exerted on protrudent parts of the copper foil surface along the direction of buff rotation, so that protrudent parts of the copper foil surface are likely to deform along the direction of buff rotation. It is difficult to effect uniform roughening treatment on the buffed copper foil having thus deformed protrudent parts. Nonuniform roughening treatment would invite problems such that the adherence to insulating films, etching uniformity, bonding reliability, etc. are deteriorated. These problems are likely to occur especially when thin electrodeposited copper foils are mechanically polished.\nAs apparent from the above, all the mechanical polishing, chemical polishing and electrolytic polishing methods, conventionally employed in the polishing of the matte side of the electrodeposited copper foil, have been highly useful for the polishing of electrodeposited copper foil for forming wiring patterns of 100 xcexcm or more pitch. However, the single employment of the above conventional polishing methods in the polishing of the matte side of the electrodeposited copper foil used to produce a printed wiring board whose pitch is becoming tens of microns or less involves the highly probable danger of causing the formed wiring pattern to be brittle and hence tends to render it difficult to supply printed wiring boards having reliable properties.\nAn object of the present invention is to provide an electrodeposited copper foil which is especially suitable for producing a printed wiring board of fine pitch and to provide a process for producing such an electrodeposited copper foil.\nAnother object of the present invention is to provide an electrodeposited copper foil with its matte side surface prepared, the preparing performed by a plurality of different polishing operations under extremely mild conditions in place of a single polishing operation to thereby realize a matte side surface prepared with high uniformity, and to provide a process for producing such an electrodeposited copper foil.\nFurther objects of the present invention are to provide a printed wiring board of stable quality produced with the use of the above electrodeposited copper foil with its surface prepared with high uniformity and to provide a multilayer laminate printed wiring board consisting of a laminate of such printed wiring boards.\nThe process for producing an electrodeposited copper foil with its surface prepared according to the present invention comprises the steps of:\nsubjecting an electrodeposited copper foil having a shiny side and a matte side whose average surface roughness (Rz) is in the range of 2.5 to 10 xcexcm to at least one mechanical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 1.5 to 3.0 xcexcm, and\nsubjecting the matte side having undergone the mechanical polishing to a selective chemical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 0.8 to 2.5 xcexcm.\nThe electrodeposited copper foil with its surface prepared according to the present invention is one having a shiny side and a prepared matte side whose average surface roughness (Rz) is in the range of 0.8 to 2.5 xcexcm, the above prepared matte side obtained through the steps of subjecting a matte side of electrodeposited copper foil whose average surface roughness (Rz) is in the range of 2.5 to 10 xcexcm to at least one mechanical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 1.5 to 3.0 xcexcm, and subjecting the matte side having undergone the mechanical polishing to a selective chemical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 0.8 to 2.5 xcexcm.\nThe prepared matte side surface of the electrodeposited copper foil with its surface prepared according to the present invention is preferably furnished with a roughened layer. A corrosion preventive layer is preferably disposed on the surface of the roughened layer, and a silane coupling agent layer is preferably disposed on the surface of the corrosion preventive layer.\nThe printed wiring board of the present invention comprises an insulating substrate having its surface furnished with a wiring pattern formed from an electrodeposited copper foil with its surface prepared, this electrodeposited copper foil with its surface prepared having a shiny side and a prepared matte side whose average surface roughness (Rz) is in the range of 0.8 to 2.5 xcexcm, the prepared matte side obtained through the steps of subjecting a matte side of electrodeposited copper foil whose average surface roughness (Rz) in the range of 2.5 to 10 xcexcm to at least one mechanical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 1.5 to 3.0 xcexcm, and subjecting the matte side having undergone the mechanical polishing to a chemical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 0.8 to 2.5 xcexcm.\nThe multilayer printed wiring board of the present invention comprises a laminate of a plurality of boards of given thickness, the above boards capable of being electrically connected to each other in the direction of the thickness of the laminate, each of the boards having its surface furnished with a wiring pattern formed from an electrodeposited copper foil with its surface prepared, this electrodeposited copper foil with its surface prepared having a shiny side and a prepared matte side whose average surface roughness (Rz) is in the range of 0.8 to 2.5 xcexcm, the prepared matte side obtained through the steps of subjecting a matte side of electrodeposited copper foil whose average surface roughness (Rz) is in the range of 2.5 to 10 xcexcm to at least one mechanical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 1.5 to 3.0 xcexcm, and subjecting the matte side having undergone the mechanical polishing to a chemical polishing so that the average surface roughness (Rz) of the matte side becomes in the range of 0.8 to 2.5 xcexcm.\nIn the present invention, it is preferable for the electrodeposited copper foil with its surface prepared that the sharp edges of the matte side have been substantially eliminated by the chemical polishing after at least one mechanical polishing.\nThe electrodeposited copper foil with its surface prepared, employed for forming the printed wiring board or the multilayer printed wiring board, preferably has its chemically polished matte side treated for roughening.\nIn the present invention, a copper-clad laminate comprises a substrate and the electrodeposited copper foil with is surface prepared which is laminated on the substrate.\nThere is a vast plurality of unevenness on the surface of the matte side of the electrodeposited copper foil, wherein, for example, protrudent part heights and depressed part depths are not fixed. The surface condition of the matte side is often expressed by the average surface roughness (Rz). As described later, the average surface roughness (Rz) is an average of depths of depressed parts and heights of protrudent parts constituting the surface of the matte side. Thus, on the matte side having a given average surface roughness (Rz), there is a multiplicity of protrudent parts whose roughness (Rmax) is greater than the average surface roughness (Rz). In the first step of the process for producing an electrodeposited copper foil with its surface prepared according to the present invention, mainly apex portions of the protrudent parts are mechanically removed by the mechanical polishing. Thus, the matte side is freed mainly of apex portions of the protrudent parts. In this mechanical polishing, stress is applied to portions to be polished so that, referring to, for example, FIG. 4(a), apex portions of the protrudent parts 112 are removed. Accordingly, there inevitably occurs a directionality along which mechanical polishing stress is applied. For example, referring to FIG. 4, when apex portions of the protrudent parts are polished by a rotary buff, stress is applied to the apex portions of the protrudent parts 112 in the same direction as the direction of buff rotation with the result that, as shown in FIG. 4(b), the protrudent part portion first brought into contact with the rotary buff is polished more easily than the side later brought into contact with the rotary buff. Furthermore, apex portions of the protrudent parts 112 are ground while being slightly deformed along the direction of stress application, so that polished upper portions of the protrudent parts may suffer from deformation corresponding to the stress exerted by the rotary buff or the like. For example, referring to FIG. 4(b), it has been found that the protrudent parts may be deformed along the direction of buff rotation with the result that shelf-shape deformations 111 are formed downstream in the direction of stress application.\nAs apparent from the above, stress is produced during the mechanical polishing, so that the mechanical polishing is inevitably accompanied by a directionality. Further, in the process of the present invention, it is intended to mainly polish and remove apex portions of the protrudent parts of the matte side by the mechanical polishing, but not to polish the whole surface of the matte side so as to render the surface specular. Rather, from the viewpoint of attaining stable bonding to the surface of an insulating substrate, it is preferred that the matte side surface be not completely smooth. Therefore, in the present invention, apex portions of the protrudent parts are mainly mechanically polished. However, referring to FIG. 4(b), this mechanical polishing forms boundary 120 between polished portion and portion not polished. When the matte side is directly subjected to the subsequent treatment for roughening, the states of being roughened may be different between the polished portion and the portion not polished.\nTherefore, in the present invention, referring to FIG. 4(c), the matte side having undergone the above mechanical polishing is further chemically polished to thereby remove discontinuous polishing boundary 120 and instead create round curved surface 121. Moreover, performing the chemical polishing after the mechanical polishing can also correct the directionality of polished surface due to stress, generated during the mechanical polishing."}
{"text": "Two interconnected PCI Express components are required to form at least one PCI Express link per port. Ports with more than a single lane may form a single link, which uses all available lanes or may optionally be sub-divided into multiple links with each link connected to a different PCI Express component. However, all traffic between any two PCI Express components travels on the same link, which can inefficiently use bandwidth and power as small data packets with low priority travel on the same link and at the same data rate as large data packets with high priority. Therefore, there is a need for a method to more efficiently utilize bandwidth and power."}
{"text": "As is known, wireless communications are becoming an increasingly indispensible aspect of contemporary life. Accordingly, systems, methods and techniques that enhance the efficiency of wireless communications systems are a welcome addition to the art."}
{"text": "Known electronic devices, generally of parallelepipedal shape, comprise a support board, including an electrical connection network, an integrated circuit chip mounted on one of the faces of the support board and an encapsulation block in which the chip is embedded. The chip is linked to the network of the support board by electrical connection elements, such as balls, interposed between the support board and the chip, or by electrical connection wires embedded in the encapsulation block.\nSuch electronic devices are mounted on printed circuit boards, generally by way of electrical connection elements, such as balls, linking the electrical connection network of the support boards and the electrical connection network of the printed circuit boards.\nWhen the chips generate radiofrequency signals having to be transmitted or when they process radiofrequency signals received, the transmit or receive antennas are produced on the printed circuit boards. The electrical signals follow very long resistive paths which consist of lines of the electrical connection network of the printed circuit boards, the electrical connection elements between the printed circuit boards and the support boards, lines of the electrical connection network of the support boards and the electrical connection elements between the support boards. Such paths depend furthermore on the quality of the interconnections arising from fabrication.\nThe arrangements hereinabove constitute a handicap notably when the necessary dimension of the antennas, for the transmission of radiofrequency signals at frequencies of the order of a gigahertz or greater, or indeed much greater, than a gigahertz, becomes reduced."}
{"text": "1. Field of the Invention\nThe present invention relates to a heat exchanger for exhaust heat recovery of a vehicle. More particularly, the present invention relates to an exhaust heat recovery apparatus for a vehicle in which a coolant and a lubricant are subjected to heat exchange for exhaust heat.\n2. Description of Related Art\nIn general, a heat exchanger in a vehicle recovers exhaust heat by performing heat exchange between an exhaust gas and a coolant in such a manner as to lower the temperature of a high-temperature exhaust gas while increasing the temperature of a low-temperature coolant.\nThe heat, recovered through the heat exchanger, is transmitted to a part of a vehicle requiring heat, and used for the heating of a lubricant, the fast warming up of an engine, heating, and the like.\nThe heat exchanger is provided with a heat exchange part in which heat exchange with a coolant takes place. An exhaust gas exchanges its heat with the coolant while flowing through the heat exchange part. In a case where heat recovery is unnecessary, the heat exchanger changes a flow path such that the exhaust gas bypasses the heat exchange part.\nFor example, as for a structure for heating a lubricant by recovering the exhaust heat, the exhaust heat is made to pass through the heat exchange part. After the lubricant is heated to an appropriate temperature, the exhaust gas is made to bypass the heat exchange part to thus prevent a further increase in the temperature of the lubricant.\nThat is, at a low temperature or in initial starting, the high-temperature exhaust heat is supplied to the heat exchange part for the exchange of heat with the coolant. The coolant having increased in temperature increases the temperature of the lubricant. Since the lubricant has a high level of viscosity at a low temperature, the lubricant increases friction when supplied to a transmission or a work part of an engine. The lubricant heated by the coolant can reduce friction of the transmission, and the like.\nThe lubricant, when excessively heated, is degenerated and thus fails to properly perform lubrication. For this reason, at a high temperature or in acceleration, the exhaust gas is made to bypass the heat exchange part. In this case since heat exchange does not occur between the exhaust gas and the coolant, the coolant is maintained at a low temperature and the lubricant can be cooled by the low-temperature coolant.\nIn such a heat exchanger according to the conventional art, a single passage through which the coolant flows is disposed at a region in which the exhaust gas passes, and a tube through which the lubricant is distributed is installed within the passage to thus form the heat exchange part.\nAccordingly, at the low temperature or in the initial starting of a vehicle, heat exchange between the exhaust gas and the coolant occurs in the region where the exhaust gas passes so as to increase the temperature of the coolant. The coolant can increase the temperature of the lubricant. The lubricant the temperature of which has been increased in the above manner contributes to enhancing fuel efficiency by reducing friction of the transmission, and the like.\nHowever, since the heat exchanger according to the conventional art has the heat exchange part formed simply by installing a tube through which the lubricant is distributed within a single coolant path, the contact area and contact time of the coolant path with respect to the exhaust gas are insufficient. This may deteriorate heat-exchange performance and slow down the temperature increase of the lubricant, resulting in deterioration of fuel efficiency.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the prior art already known to a person skilled in the art."}
{"text": "Technical Field\nThe current disclosure relates generally to a new embodiment of bulk GaAs with an increased carrier lifetime of at least 10 microseconds to be used for optical and electrical devices. The current disclosure also relates to the apparatus, systems, and methods for creating GaAs with very long free carrier lifetimes. More particularly, this very long free carrier lifetime GaAs is expected to be useful as a semiconductor radiation detector material, and also is expected to be useful for applications including, but not limited to, medical imaging, solar cells, diode lasers, and optical limiters.\nBackground Information\nGaAs is a well-known semiconductor and is grown by many methods. It can be produced using a variety of techniques including both bulk melt growth and vapor growth. Commercially available GaAs always has a significant concentration of a defect called EL2 (As on Ga antisite defects) which are known to greatly reduce the free-carrier lifetime of the material. The EL2 antisite defects form deep level traps and are inherent in any melt grown material due to a widening of the solubility curve of As in GaAs as the GaAs crystal temperature is near that of the melting point of the compound.\nCrystal growth from the vapor phase can be done at a temperature significantly lower than that of the melting point of the compound, however, most vapor growth techniques such as metal organic chemical vapor deposition (MOCVD) and molecular beam epitaxy (MBE) have growth rates that are much too slow (approximately 1 micron per hour) to make growth of any material approaching that of bulk quantities (>100 micrometers) practical.\nGaAs with lower concentrations of EL2 defects or longer free carrier lifetimes would be useful because the electrons generated in a process like the absorption of radiation would be able to travel further distances through the GaAs before being trapped by a defect.\nOne area where longer lifetime (low EL2 defect) bulk GaAs is expected to be useful is in the manufacture of semiconductor radiation detectors. Improving the energy resolution of gamma radiation detectors is among the most important science and technology objectives for national security applications since it enables the use of high-resolution energy spectroscopy to distinguish between the natural radioactive isotopes, medical or commercially used radioisotopes, and radioisotopes that pose a threat. Gamma radiation detectors must be very thick—on the order of several millimeters—because of the deep penetration of high energy gamma rays which must be fully absorbed for accurate quantitative detection.\nOf the two primary classes of radiation detector materials—semiconductors and scintillators—semiconductors are fundamentally capable of much higher energy resolution because the information carriers are the electrons and holes directly produced by the energy cascade. Desirable properties in such a semiconductor are a room temperature operation—which requires a band gap between 1.35-2.55 eV (McGregor and Hermon, Nuclear Instruments and Methods in Physics Research Section A: Accelerators, Spectrometers, Detectors and Associated Equipment Volume 395, Issue 1, Room-temperature compound semiconductor radiation detectors, pages 101-124 (1997), high electron and hole mobilities (μ), and long carrier lifetimes (τ). Historically, however, semiconductor radiation detectors have been limited by high cost, operational complexity, low efficiency due to limited size, and difficulty achieving high-purity, defect-free crystals. Unfortunately, progress in overcoming these obstacles over the years has been slow and incremental.\nProgress in the development of high resolution, room temperature, gamma radiation detectors has been severely limited for many years by the size, quality, and transport properties of single crystal compound semiconductors. CdxZn1-xTe (CZT) is the most commercially advanced of these materials, but Te inclusions, twins, and grain boundaries are constant barriers to the size, yield, and cost of these crystals for detector applications.\nGallium arsenide (GaAs) has been studied as a radiation detector since the early 1960's, predating CZT, but its widespread use for gamma ray detection has never been realized due to the presence of native deep level traps (EL2, i.e., AsGa antisites) which reduced its free carrier lifetime at room temperature. Otherwise, GaAs has very attractive intrinsic properties. Its band gap of 1.42 electron voltage (eV) is near optimum. Theory predicts that carrier lifetimes in very pure and well-ordered GaAs should approach 100 microseconds or more, but such long lifetimes have never historically been observed.\nA need exists, therefore, for bulk GaAs with free-carrier lifetimes significantly greater than that of the previous state of the art of 1 microsecond."}
{"text": "(a) Field of the Invention\nThis invention relates to a novel process for the preparation of [bis(substituted aryl)(indolyl)]methanes and [(substituted aryl)(heteryl)(indolyl)]methanes useful in the art of carbonless duplicating as, for example, in pressure-sensitive systems and in thermal marking systems; and to novel [(aryl)(disubstituted aminoaryl)(indolyl)]methanes, [(aryl)(disubstituted aminoaryl)(indolyl carboxylic acid)]methanes and [(aryl)(julolidino)(indolyl)]methanes produced by said process.\n(b) Description of the Prior Art\nSeveral classes of organic compounds of widely diverse structural types are known to be useful as colorless precursors for carbonless duplicating systems. Among the more important classes, there may be named phenothiazines, for example, benzoyl leuco methylene blue; phthalides, for example, crystal violet lactone; fluorans, for example, 2'-anilino-6'-diethylaminofluoran and 2'-dibenzylamino-6'-diethylaminofluoran; arylsulfinate salts of Michler's Hydrol; and various other types of colorless precursors currently employed in commercially-accepted carbonless copy systems. Typical of the many such systems taught in the prior art are those described in U.S. Pat. Nos. 2,712,507, 2,800,457, 3,041,289 and 4,000,087, which issued July 5, 1955, July 23, 1957, June 26, 1962 and Dec. 28, 1976, respectively. Many of the color formers in the prior art suffer one or more disadvantages such as low tinctorial strength, poor light stability, low resistance to sublimation, low susceptibility to copiability of the color-developed form in standard office copying machines, for example, a xerographic type of copier, and low solubility in common organic solvents, the latter disadvantage thus requiring the use of specialized and expensive solvents in order to obtain microencapsulated solutions of sufficient concentration for use in pressure-sensitive copying systems.\nIn the prior art, benzoyl leuco methylene blue, related derivatives of phenothiazines, and analogous compounds derived from phenothiazines have been used as slow developing color formers for use in carbonless duplicating systems in combination with additional art recognized instant color-forming compounds.\nOne of the oldest combinations of color-forming compounds used in pressure-sensitive systems is a mixture of crystal violet lactone and benzoyl leuco methylene blue. The crystal violet lactone is responsible for providing the initial image which is, however, very unstable to light and moisture. These deficiencies are overcome by the use of benzoyl leuco methylene blue which develops slowly upon exposure to air, providing a green-blue image to augment or replace that of the fading crystal violet lactone. The developed image of the benzoyl leuco methylene blue is extremely light stable, but is rather green in shade, is lacking in contrast quality, and therefore is unsatisfactory when multiple copies are required. Further, it does not possess adequate xerographic copiability on the commercially-available reproduction machines.\nThe following items to date appear to constitute the most relevant prior art with regard to the instant invention.\nU.S. Pat. No. 3,637,748, issued Jan. 25, 1972, discloses and claims a heteryl methane compound which is free of sulphonic acid groups and has the formula ##STR1## wherein R.sub.1 is hydrogen or methyl; R.sub.2 is methyl or phenyl; B stands for a member selected from the group consisting of phenyl, nitrophenyl, chlorophenyl, dichlorophenyl, methoxyphenyl, methylphenyl, hydroxyphenyl, thienyl-2, and furyl-2; D stands for a member selected from the group consisting of 4-methylquinolinyl-methosulphate, 1,2,3,3-tetramethyldihydroindolyl, 1-phenyl-3-methylpyrazolone-(5)-yl, 1-phenyl-2,3-dimethyl-pyrazolone-(5)-yl, substituted naphthyl, substituted phenyl, substituted chlorophenyl, substituted methylphenyl, substituted hydroxyphenyl and substituted ethoxyphenyl, wherein the substituent is located in the 4-position relative to the position at which the CH-group is bound and is a member selected from the group consisting of phenylamino, ethoxyphenyl, amino, methylphenylamino, methylethoxyphenylamino, piperidino, morpholino, mono-alkylamino wherein the alkyl portion contains 1-4 carbon atoms, di-lower alkylamino wherein the lower alkyl portion contains 1-2 carbon atoms, and said mono-alkylamino and said di-lower alkylamino substituted by a member of the group consisting of hydroxy, chloro, cyano and dimethylamino. The heteryl methane compounds are disclosed as being valuable dyestuff intermediates.\nU.S. Pat. No. 3,642,823, issued Feb. 15, 1972, discloses and claims as aminodiphenyl-indolyl-methane dyestuff free of sulfonic acid and carboxylic acid groups having the formula ##STR2## wherein R is hydrogen, lower alkyl containing 1-5 carbon atoms, nitro, cyano, benzyl, phenyl, carboxylic acid methyl ester, carboxylic acid ethyl ester, unsubstituted carbonamido, substituted carbonamido containing N-substituents selected from the group consisting of lower alkyl, and lower alkyl substituted with Cl, Br, or OH, methylsulfonyl, ethylsulfonyl, phenylsulfonyl, p-toluene sulfonyl, benzoyl, or acetyl; R.sub.1 and R.sub.2 are both hydrogen; or R.sub.1 is unsubstituted lower alkyl containing 1-5 carbon atoms, substituted lower alkyl containing 1-5 carbon atoms and substituted with Cl, Br or OH, benzyl, naphthyl, phenyl or phenyl substituted with Cl, Br, lower alkyl or lower alkoxy and R.sub.2 is unsubstituted lower alkyl containing 1-5 carbon atoms, substituted lower alkyl containing 1-5 carbon atoms and substituted with Cl, Br, OH, or CN or benzyl; R.sub.3 is hydrogen, lower alkyl containing 1-5 carbon atoms, lower alkoxy, phenyl, phenyl substituted with chlorine, bromine, lower alkyl or lower alkoxy or carboxylic acid methyl or ethyl ester; R.sub.4 is hydrogen, lower alkyl containing 1-5 carbon atoms, lower alkyl containing 1-5 carbon atoms substituted with CN, or benzyl; K, L and M are independently selected from the group consisting of hydrogen, lower alkyl, lower alkoxy, nitro, Cl, and Br, or two of the substituents M are fused to form a naphthalene nucleus; p, q and r are 1 or 2; and X is a substantially colorless anion which does not impair the solubility of the dyestuff in an undesirable manner. The aminodiphenyl-indolyl-methane dyestuffs are described as being useful in the dyeing and printing of materials of leather, tannin-treated cotton, cellulose acetate, synthetic superpolyamides and superurethanes as well as for the dyeing of lignin-containing fibres, such as coconut fibres, jute and sisal. They are also suitable for the production of liquid writing inks, stamping inks, ball point pen pastes and they may also be used in offset printing. Above all they are eminently suitable for the dyeing and printing of materials which are composed partially or completely of polymerized acrylonitrile and/or vinylidene cyanide or which are composed partially or completely of acid-modified aromatic polyesters, such as sulphonic acid group-containing polyethylene terephthalate.\nU.S. Pat. No. 3,995,088, issued Nov. 30, 1976, discloses and claims a pressure-sensitive recording material comprising at least one pair of paper sheets containing at least two color formers, dissolved in an organic solvent, and an electron-accepting substance wherein at least one color former is a leuco methylene dyestuff of the formula ##STR3## wherein Y represents an amino-substituted phenyl residue of the formula ##STR4## or an indolyl residue of the formula ##STR5## wherein R.sub.1 and R.sub.2, independently of the other represent hydrogen, alkyl with 1-12 carbon atoms, alkoxyalkyl with 2-8 carbon atoms, benzyl or phenyl; R.sub.3 is hydrogen, halogen, nitro, alkyl having 1-4 carbon atoms or alkoxy having 1-4 carbon atoms; X.sub.1 represents hydrogen, alkyl with 1-12 carbon atoms, alkenyl with at most 12 carbon atoms or benzyl; X.sub.2 represents hydrogen, alkyl with 1-12 carbon atoms or phenyl; the ring A is unsubstituted or substituted by cyano, nitro, halogen, alkyl having 1-4 carbon atoms, alkoxy having 1-4 carbon atoms or acyl with 1-8 carbon atoms; and Z represents alkyl with 1-12 carbon atoms, alkenyl with at most 12 carbon atoms, aryl, aralkyl, a heterocyclic radical or the residue of an organic, particularly aliphatic or cycloaliphatic, compound having a ketomethylene group. The other color former is selected from the group consisting of crystal violet lactone, benzoyl leuco methylene blue, diamino substituted fluoran compounds, 3-phenyl-3-indolylphthalides or 3,3-bis-indolylphthalides or their mixtures thereof.\nWhitehead and Whitesilt in the Journal of Medicinal Chemistry 1974 Vol. 17, No. 12, 1298-1304 (1974) describe the preparation and physical characteristics of [bis(phenyl)(2-carboxyindol-3-yl)]methane from the interaction of 2-carboxyindole and benzhydrol in glacial acetic acid. The compound and similar compounds were screened for antimicrobial activity. Toxicity, behavioral effects, antiviral activity, antiinflammatory properties, antispasmatic activity and general endocrine properties of the (diarylmethyl)indoles are described.\n(c) Prior Publication\nWest German Offenlegungsschrift No. 2,917,271, which was published on Nov. 6, 1980, describes pressure- and heat-sensitive copy material, characterized by the fact that it contains a color component of the formula ##STR6## in which R.sup.1 designates hydrogen, hydroxy, C.sub.1 to C.sub.6 alkoxy, C.sub.2 to C.sub.6 alkenyloxy, cycloheyloxy, cyclopentyloxy, possibly benzyloxy, phenylethyloxy, phenyloxy, naphthyloxy, phenylthio or naphthylthio, C.sub.1 to C.sub.6 alkylthio substituted by C.sub.1 to C.sub.4 alkyl, C.sub.1 to C.sub.4 alkoxy or halogen; R.sup.2 and R.sup.3 independently of one another designate hydrogen, or C.sub.1 or C.sub.6 alkyl possibly substituted by C.sub.1 to C.sub.4 alkoxy, C.sub.1 to C.sub.4 alkyl, carbonyloxy, hydroxy, halo or cyano, or benzyl, phenylethyl, phenyl or naphthyl, cyclopentyl or cyclohexyl, substituted by C.sub.1 to C.sub.4 alkyl, C.sub.1 to C.sub.4 alkoxy or halogen, or R.sup.2 and R.sup.3 along with the nitrogen atom designate as pyrrolidine, pyrazoline, piperidine or morpholine ring, or R.sup.2 along with the nitrogen atom and the ring B designate an indoline or tetrahydroquinoline ring, wherein the rings mentioned can be substituted by methyl or phenyl; R.sup.4 designates C.sub.1 to C.sub.4 alkyl or phenyl; R.sup.5 designates hydrogen, possibly C.sub.1 to C.sub.22 alkyl, C.sub.2 to C.sub.4 alkenyl, cyclohexyl substituted by cyano, chloro or hydroxy, phenyl or benzyl substituted by methyl, ethyl, methoxy, ethoxy or chlorine; R.sup.6, R.sup.7 and R.sup.8 independently of one another designate hydrogen, methyl, ethyl, methoxy, ethoxy or chlorine. This reference appeared subsequent to applicants' invention described herein and less than one year prior to the filing date of this application."}
{"text": "A commonly used lawncare apparatus, such as a walk-behind lawnmower, includes a cutting device, a deck supporting the cutting device, a chassis having four wheels to which the deck is mounted, and a handle for pushing a main body of the lawncare apparatus disposed at one end of the chassis.\nWhen using the lawnmower, there are two ways to transport the lawnmower from a storage room to a lawn. The first way is to place the four wheels on the ground and push/pull the lawnmower. The second way is to lift the lawnmower away from the ground. The first way requires a large storage room and the second way requires that the lawnmower is of light weight. That is, the way to transport the lawnmower from the storage room to the lawn needs to be improved."}
{"text": "Modern computing systems provide various methods for user authentication. Some authentication methods include a password requirement, where the user enters his or her user identifier as well as a secret password. This is referred to as single-factor authentication. More secure authentication regimes, such as multiple-factor authentication, require, in addition to a password, a second means of verification, such as a token, smartcard, fingerprint, retinal scan, etc. For example, with smartcard authentication, a user inserts a smartcard into a smartcard reader and enters a personal identification number (PIN). When a correct PIN is entered, one or more certificates that are stored in the smartcard are used to authenticate the user.\nIn computing systems that access remote resources, a user may need to be authenticated each time the remote resources are accessed, requiring the user to repeatedly enter his or her credentials. As a result, for user convenience, a feature referred to as single sign-on (SSO) may be used to reduce the burden on the user by enabling user credentials to be shared between and reused by multiple applications. For example, a web browser that uses SSO may rely on a JSON Web Token (JWT) that identifies the user and stores an authentication event. The JWT can then be stored in a browser cookie for subsequent use by other applications that require a similar type of authentication. That is, a request to the authentication service on behalf of other applications executed within the same browser may include the browser cookie, enabling the authentication service to validate the JWT and verify that the user has already authenticated.\nIn enterprise applications, SSO is an increasingly popular feature due to trends such as Bring Your Own Device (BYOD) and Corporate-Owned Personally Enabled (COPE) devices. For example, enterprise administrators may require enforcement of specific authentication policies when a user accesses sensitive resources, such as payroll, banking, and health care. Consequently, SSO provides a convenient way for a user to enter his or her credentials once and reuse the stored credentials to access various types of sensitive resources.\nHowever, because native mobile applications do not store tokens in a browser, the native applications cannot share authentication information in a cookie associated with a shared authentication service. Accordingly, one common mobile SSO approach is token sharing. Token sharing occurs when one native application obtains a token from an authentication service (AS) and then shares that token with other mobile applications. Because the token is shared, the token works only in the situation where all applications are owned by one enterprise resource server. Additionally, because the AS does not know which application submitted the token to access the resource server, the AS cannot revoke the token to prevent only that application from accessing resources, since doing so would block all applications.\nFurther, in some circumstances, even if the user has already been authenticated to the SSO system, accessing sensitive resources may require further authentication via a more secure authentication method, commonly referred to as step-up authentication. As a result, the user is still required to enter his or her credentials each time authentication via the more secure method is requested by an application. Moreover, after the user has provided his or her credentials to perform step-up authentication for a particular application, the user may have to re-enter those credentials (e.g., re-authenticate) to access a resource with that application each time the user's session expires.\nAccordingly, there are challenges faced by system administrators when attempting to provide users with convenient and secure access to sensitive resources."}
{"text": "The present invention relates to a method of determining the concentration of heparin in fluid samples, and to a device for performing the method.\nHeparin is a heteropolysaccharide of the acid mucopolysaccharide type which has anticoagulant activity resulting from its ability to catalyze the reaction between antithrombin III and thrombin. On the basis of this activity, heparin is widely used as an anticoagulant in cardiovascular surgery, for example, during operations with extracorporeal circulation, and in other diagnostic and therapeutic applications.\nWhen heparin is used as an anticoagulant, it may be very useful to have available a method which permits rapid, reproducible and accurate measurement of the concentration of heparin present in the sample of interest, and a device for performing this measurement.\nIn order to perform this measurement, it is known to use dyes of the cationic thiazines series. As a result of the reaction of these dyes with heparin, it is possible to observe a change in the absorption (or transmission) spectrum of the dye, that is, a reduction in the absorption due to the free dye in solution and the appearance of an absorption band of the heparin-dye complex formed. The absorption spectrum of cationic thiazines in dilute aqueous solution is in fact characterized by a main absorption band due to the contribution of the monomer and the dimer of the dye in solution. When heparin is added to this dilute solution, a second band (called the xcexc band) which corresponds to the heparin-dye complex formed appears in the absorption spectrum of the solution. The addition of heparin to the dye solution not only causes the appearance of the xcexc band, but also causes a reduction in the absorption of the main band. The reduction in the absorption of the free dye in solution and the increase in the absorption of the dye-heparin complex constitute the two components of a phenomenon known by the term xe2x80x9cmetachromasiaxe2x80x9d, the magnitudes of which, in suitable reaction conditions, can be correlated with the concentration of heparin present in the sample.\nOn the basis of this principle, assays have been developed for measuring the concentration of heparin in samples of interest. U.S. Pat. No. 4,911,549 (Karkar) describes a method of determining the heparin concentration in blood plasma by metachromatic reaction with the dye Azure A which belongs to the cationic thiazines series. This assay is based on the measurement of two distinct transmittance signals for two distinct wavelengths, one of which is substantially insensitive to the heparin dilution.\nIt has surprisingly been found that the measurement of the entire absorption spectrum of the dye in a metachromatic heparin assay, rather than measurement of distinct absorption signals with predetermined separate wavelengths, achieves results which are more accurate and reproducible and, above all, are independent of the reaction medium. In the case of assays performed on samples of biological fluids such as blood plasma containing proteins and other chemical species which could potentially interfere with accurate measurement of the absorption at these predetermined wavelengths, this latter characteristic is of substantial importance because it makes possible a device which can perform the assay automatically and which does not require repeated calibrations.\nAccordingly, the invention provides a method for determining the concentration of heparin in a fluid sample comprising: (a) providing a fluid sample containing heparin; (b) adding to the fluid sample a solution of a dye to form a mixture of sample and dye, wherein the dye interacts with the heparin in the sample so that the absorption spectrum of the mixture of sample and dye in the visible range varies as a result of the interaction in a manner quantitatively dependent on the heparin concentration; (c) determining the absorption spectrum in the visible range of the mixture of sample and dye; and (d) calculating a spectral parameter representative of both the reduction in the absorption of the free dye in solution and the increase in the absorption of the dye-heparin complex, the value of the spectral parameter being indicative of the concentration of heparin present in the fluid sample, in order to determine the concentration of heparin present in the fluid sample. The relationship between the value of the spectral parameter and the concentration of the heparin present in the fluid sample has been previously determined by determining the absorption spectra, in the visible range, of a composition comprising the fluid and the dye, in the absence of heparin and in the presence of a plurality of concentrations of heparin, and calculating the relationship between the value of the spectral parameter and the concentration of heparin. A device for performing this method is also provided.\nAdditional features and advantages of the invention are set forth in the description which follows and in part will be apparent from the description. The objectives and other advantages of the invention will be realized and attained by the method of determining the concentration of heparin in a fluid sample and the device for performing this method as particularly pointed out in the written description and claims.\nIt is to be understood that both the foregoing general description and the following detailed description are exemplary and explanatory and are intended to provide further explanation of the invention as claimed."}
{"text": "Tobacco (Nicotiana tabacum L.) is an important commercial crop in the United States as well as in other countries. In tobacco plants, N-demethylation of nicotine results in nornicotine, a secondary alkaloid known to be a precursor for formation of N-Nitrosonornicotine (“NNN”) in cured leaves. At high levels, NNN is an undesired component of cured leaves.\nThe predominant alkaloid found in commercial tobacco varieties is nicotine, typically accounting for 90-95% of the total alkaloid pool. The remaining alkaloid fraction is comprised primarily of three additional pyridine alkaloids: nornicotine, anabasine, and anatabine. Nornicotine is generated directly from nicotine through the activity of the enzyme nicotine N-demethylase. Nornicotine usually represents less than 5% of the total pyridine alkaloid pool, but through a process termed “conversion,” tobacco plants that initially produce very low amounts of nornicotine give rise to progeny that metabolically “convert” a large percentage of leaf nicotine to nornicotine. In tobacco plants that have genetically converted (termed “converters”), the great majority of nornicotine production occurs during the senescence and curing of the mature leaf (Wernsman and Matzinger (1968), Tob. Sci., 12:226-228). Burley tobaccos are particularly prone to genetic conversion, with rates as high as 20% per generation observed in some cultivars.\nDuring the curing and processing of the tobacco leaf, a portion of the nornicotine is metabolized to the compound NNN, a tobacco-specific nitrosamine (TSNA) that has been reported to be carcinogenic in laboratory animals (Hecht and Hoffmann (1990), Cancer Surveys, 8:273-294; Hoffmann et al. (1994), J. Toxicol. Environ. Health, 41:1-52; Hecht (1998), Chem. Res. Toxicol., 11:559-603). In flue-cured tobaccos, TSNAs are found to be predominantly formed through the reaction of alkaloids with the minute amounts of nitrogen oxides present in combustion gases formed by the direct-fired heating systems found in traditional curing barns (Peele and Gentry (1999), “Formation of Tobacco-specific Nitrosamines in Flue-cured Tobacco,” CORESTA Meeting, Agro-Phyto Groups, Suzhou, China). Retrofitting these curing barns with heat-exchangers virtually eliminated the mixing of combustion gases with the curing air and dramatically reduced the formation of TSNAs in tobaccos cured in this manner (Boyette and Hamm (2001), Rec. Adv. Tob. Sci., 27:17-22.). In contrast, in the air-cured Burley tobaccos, TSNA formation proceeds primarily through reaction of tobacco alkaloids with nitrite, a process catalyzed by leaf-borne microbes (Bush et al. (2001), Rec. Adv. Tob. Sci., 27:23-46). Thus far, attempts to reduce TSNAs through modification of curing conditions while maintaining acceptable quality standards continue to be challenging for the air-cured tobaccos."}
{"text": "The present invention relates to cotton harvesters and, in particular, to a scrapping plate assembly for a cotton harvester row unit.\nCotton harvesting has always been economically sensitive. When productivity increases were necessary to meet demand, manual picking was replaced by machinery. Though requiring less labor content, the machinery was capital intensive and less efficient in harvesting the available crop, failing to pick immature cotton bolls and dropping to the ground otherwise utile product. Gains were made in production speed and harvest efficiency as multi-row harvesters became available. Notwithstanding these improvements, to remain competitive in cost and supply, dual pass harvesting was employed to reap the cotton missed in the first harvest.\nAs the cotton sector evolves, smaller tracts of land are being supplanted by large fields, typically 1,000 acres or more. When all the costs and times associated with the larger fields have been considered, it is generally deemed uneconomical to conduct multiple pass harvesting, and the dropage and incomplete harvesting considered inescapable yield losses. Experience has shown that as much as 10% to 15% of the cotton lint remains on the stalks or on the ground after a single harvesting path.\nAccessories such as scrapping plates have been employed for increasing efficiencies, but have not been considered acceptable for initial harvesting particularly for heavy cotton or use with short knotty cotton. These scrapping plates carried on the pressure plates are used in conjunction with spindle-type pickers, such as the John Deere 9965 series and as disclosed in U.S. Pat. Nos. 5,203,152 to Deutsch et al.; U.S. Pat. No. 5,519,988 to Copley; U.S. Pat. No. 5,557,910 to Del Rosario; U.S. Pat. No. 6,293,078 to Deutsch el al.; and U.S. Pat. No. Re. 32,699 to Fachini. Therein, vertical rows of rotating spindles carried on an orbiting drum engage the cotton bolls with pointed barbs, carry the entangled cotton between grid bars to doffer columns whereat the cotton is removed from the spindles and discharged to a vacuum conveyor system for delivery to the storage container. The scrapping plates have projecting triangular ribs that coact with alternate spindles to compressively force the cotton through the passage therebetween, engaging and capturing bolls that would otherwise not be harvested. The contact path is relatively short, about 3 inches, and shallow, about 1 inch, such that only about two barbs contact the compressed boll thereby limiting engagement with the result that cotton bolls may drop to the ground or be dislodged by the grid bars.\nIn view of the foregoing, it would be desirable to provide a cotton harvester adjunct for use in single pass cotton harvesting to provide increased yields through reduction in dropage and an increase in stalk picking.\nThe foregoing is accomplished in the present invention by a scrapping plate assembly that employs elongated tapered ribs adjacent each passing spindle to increase the contact path while the cotton in a compacted condition thereby increasing barb engagement and contact time with the spindle, increasing harvesting from the stalks and reducing dropage. The scrapping plate assembly includes a base plate, fixed at existing mounting holes on the rear of the row unit pressure plate, and carrying a plurality of elongated ribs providing independent parallel paths for each spindle. The ribs have a tapered entry section for progressively engaging the product, a center section for engaging multiple barbs with the product over a substantial distance and height, and a tapered exit section extending beyond the pressure plate and gradually reducing contact with the product adjacent the grid bars. The scrapping plate assembly markedly increases the cotton removed from the stalks in a single pass and decreases the product lost through dropage.\nAccordingly, it is an object of the present invention to provide an accessory for increasing the cotton harvesting on a row unit cotton harvester.\nAnother object of the invention is to provide an improved scrapping plate yielding an improved cotton removal and retrieval in a single harvesting operation.\nA further object of the invention is to provide a scrapping assembly for a cotton harvester row unit that increases the capture of cotton by the picker spindles."}
{"text": "Programming a universal remote control may involve an overwhelming number of steps. For example, a user may be required to know a significant number of details about the device to be controlled, such as the specific device manufacturer and the specific model. In another example, the user may be required to point the original remote control at the universal remote control and press a button on the original remote control in order for the universal remote to receive and learn the infrared signal for that button. In yet another example, the user may be required to press the power button on the remote control until the device to be controlled is turned off while the remote control cycles through various output signals associated different devices. Even further, this programming process has to be repeated for each device to be controlled by the universal remote control. As the number of devices to be controlled continues to increase (e.g., a television, a set-top box, a DVD player, a digital video recording device, an audio receiver, a home theater system, etc.), programing a universal remote control is often a time-consuming experience.\nAccordingly, it is desirable to provide methods, systems, and media for providing a remote control interface."}
{"text": "1. Field of the Invention\nThe present invention relates to an automatic transmission primarily intended for motor vehicle, and more particularly, to a control valve for controlling the pressure of fluid to a torque converter in an automatic transmission.\n2. Description of Related Art\nGenerally speaking, land vehicles require three basic components. These components comprise a power plant (such as an internal combustion engine) a power train and wheels. The internal combustion engine produces force by the conversion of the chemical energy in a liquid fuel into the mechanical energy of motion (kinetic energy). The function of the power train is to transmit this resultant force to the wheels to provide movement of the vehicle.\nThe power train's main component is typically referred to as the \"transmission\". Engine torque and speed are converted in the transmission in accordance with the tractive-power demand of the vehicle. The vehicle's transmission is also capable of controlling the direction of rotation being applied to the wheels, so that the vehicle may be driven both forward and backward.\nA conventional transmission includes a hydrodynamic torque converter to transfer engine torque from the engine crankshaft to a rotatable input member of the transmission through fluid-flow forces. The transmission also includes frictional units which couple the rotating input member to one or more members of a planetary gearset. Other frictional units, typically referred to as brakes, hold members of the planetary gearset stationary during flow of power. These frictional units are usually brake clutch assemblies or band brakes. The drive clutch assemblies can couple the rotating input member of the transmission to the desired elements of the planetary gearsets, while the brakes hold elements of these gearsets stationary. Such transmission systems also typically provide for one or more planetary gearsets in order to provide various ratios of torque and to ensure that the available torque and the respective tractive power demand are matched to each other.\nTransmissions are generally referred to as manually actuated or automatic transmissions. Manual transmissions generally include mechanical mechanisms for coupling rotating gears to produce different ratio outputs to the drive wheels.\nAutomatic transmissions are designed to take automatic control of the frictional units, gear ratio selection and gear shifting. A thorough description of general automatic transmission design principals may be found in \"Fundamentals of Automatic Transmissions and Transaxles,\" Chrysler Corporation Training Manual No. TM-508A. Additional descriptions of automatic transmissions may be found in U.S. Pat. No. 3,631,744, entitled \"Hydromatic Transmission,\" issued Jan. 4, 1972 to Blomquist, et al., and U.S. Pat. No. 4,289,048, entitled \"Lock-up System for Torque Converter,\" issued on Sep. 15, 1981 to Mikel, et al. Each of these patents is hereby incorporated by reference.\nIn general, the major components featured in such an automatic transmission are: a a torque converter as above-mentioned; fluid pressure-operated multi-plate drive or brake clutches and/or brake bands which are connected to the individual elements of the planetary gearsets in order to perform gear shifts without interrupting the tractive power; one-way clutches in conjunction with the frictional units for optimization of power shifts; and transmission controls such as valves for applying and releasing elements to shift the gears (instant of shifting), for enabling power shifting, and for choosing the proper gear (shift point control), dependent on shift-program selection by the driver (selector lever), accelerator position, the engine condition and vehicle speed.\nThe control system of the automatic transmission is typically hydraulically operated through the use of several valves to direct and regulate the supply of pressure. This hydraulic pressure control will cause either the actuation or deactuation of the respective frictional units for effecting gear changes in the transmission. The valves used in the hydraulic control circuit typically comprise spring-biased spool valves, spring-biased accumulators and ball check valves. Since many of these valves rely upon springs to provide a predetermined amount of force, it will be appreciated that each transmission design represents a finely tuned arrangement of interdependent valve components. While this type of transmission control system has worked well over the years, it does have its limitations. For example, such hydraulically controlled transmissions are generally limited to one or a very small number of engines and vehicle designs. Therefore, considerable cost is incurred by an automobile manufacturer to design, test, build, inventory and repair several different transmission units in order to provide an acceptable broad model line for consumers.\nAdditionally, it should be appreciated that such hydraulically controlled transmission systems cannot readily adjust themselves in the field to compensate for varying conditions such as normal wear on the components, temperature swings and changes in engine performance over time. While each transmission is designed to operate most efficiently within certain specific tolerances, typical hydraulic control systems are incapable of taking self-corrective action on their own to maintain operation of the transmission at peak efficiency.\nHowever, in recent years, a more advanced form of transmission control system has been proposed, which would offer the possibility of enabling the transmission to adapt itself to changing conditions. In this regard, U.S. Pat. No. 3,956,947, issued on May 18, 1976 to Leising, et al., which is hereby incorporated by reference, sets forth a fundamental development in this field. Specifically, this patent discloses an automatic transmission design which features an \"adaptive\" control system that includes electrically operated solenoid-actuated valves for controlling certain fluid pressures. In accordance with this electric/hydraulic control system, the automatic transmission would be \"responsive\" to an acceleration factor for controlling the output torque of the transmission during a shift from one ratio of rotation (between the input and output shafts of the transmission) to another. Specifically, the operation of the solenoid-actuated valves would cause a rotational speed versus time curve of a sensed rotational compound of the transmission to substantially follow along a predetermined path during shifting."}
{"text": "In general, current cholangiographic agents used as biliary X-ray contrast agents are now considered too toxic to meet risk/benefit criteria for routine use as contrast agents to visualize the biliary system. Because of the relative toxicity of these agents, these agents are being replaced in use by alternative imaging methodologies, such as ultrasound and magnetic resonance imaging (MRI). On the other hand, the resolution achievable with X-ray radiology techniques is superb. Accordingly, biliary contrast agents can still be quite useful.\nA second issue with regard to these agents is the low rate of choleresis associated with biliary excretion of the agent. The low rate of choleresis increases the biliary concentration and thus increases the toxicity.\nCurrent biliary X-ray contrast agents have amphophilic properties. These agents have a hydrophilic head portion; that is, the --CO.sub.2 H terminus, and a hydrophobic tail; that is, triiodo aniline moiety with an unsubstituted 5-position. Because of this property, they are excreted nearly exclusively by the biliary route. Typical of most amphilphilic molecules, these biliary contrast agents tend to self associate. This property is reflected by the low critical micellar concentration and low osmolalities of their aqueous salt solution.\nIt is therefore desirable to produce a less toxic and more choleretic contrast agent."}
{"text": "A. General Principles\nESWL is a non-invasive technique for destroying biliary and renal concretions, i.e., stones, using acoustic shock waves. The shock waves are generated externally to the patient and focused on the site of the stone. During ESWL treatment the stones are fragmented into pieces small enough to either pass out of the body through normal excretory channels (ureter or bile ducts) or respond favorably to chemical dissolution treatment.\nThe use of sound waves for ultrasound and electrohydraulic lithotripsy by direct contact with the target was well established prior to ESWL. However, the use of shock waves in ESWL is different. Ultrasound consists of sinusoidal waves of defined wave length, with alternating positive and negative deflections. Shock waves consist of a single positive pressure front of multiple frequencies with a steep onset and a gradual decline.\nShock waves undergo less attenuation than ultrasound waves when propagated through water or body tissue. As a result, shock waves can be transmitted through water and into the body with little loss of energy or damage to tissue.\nThe use of shock waves in the medical field for the destruction of urinary stones is based on the following properties:\n1. Shock waves give rise to mechanical stress in brittle materials, such as human kidney stones. PA1 2. Shock waves lead to disintegration of such brittle material. PA1 3. Shock waves generated by the underwater discharge of a capacitor can be reliably reproduced. PA1 4. Shock wave energy can be propagated through water bath and body tissue to the stone with minimal energy loss or damage to tissue. PA1 5. Shock waves can be precisely focused by integrating the energy source with a suitable reflecting system. PA1 1. Each signal is windowed, i.e., chopped, around the arrival time of any reflected pressure waves. The duration of this window is such that the windowed signal comprises the prominent features of the reflected shock wave. PA1 2. The average amplitude of the data points comprising each signal is subtracted from each data point in that signal. In so doing, each signal is shifted such that the average value of the data points in that signal is zero. This operation is equivalent to removing any DC-shift in the signal. PA1 3. The amplitudes of the data points comprising each signal are divided by the signal's maximum positive amplitude within the chosen window. (The maximum positive value of the output from all of the sensors will be identical after this operation.) This normalization scheme facilitates the subsequent analyses. PA1 4. The windowed signal from each sensor is time-shifted with respect to the signals from each of the other sensors, in turn, in such a way that the correlation between the signals is maximized for each combination of sensors. The time shift for maximum correlation is equal to the difference in the time of arrival of the reflected pressure waves at the two sensors. These differences in arrival times (delays in arrival at one sensor with respect to the other) are then stored in the computer's memory for each sensor pairing. Thus, for a system of four pressure sensors, six delay times are calculated. PA1 5. The delay times for each sensor pairing are used to determine the coordinate components of a vector from the geometric focus of the lithotripter to the stone location.\nBased on these principles, repeated shock wave stress will eventually exceed the comprehensive strength of the stone and lead to its disintegration.\nOnce the focused shock wave reaches the stone, the pressure front is partially reflected at the front surface of the stone, thus producing compressive and tensile components, which leads to buildup of a high-pressure gradient and causes disintegration of the stone's front surface. A portion of the wave continues through the stone and is reflected at the rear surface, where the same effect takes place. The disintegration of the outer layers exposes new surfaces that in turn are broken into fragments. This process eventually results in the complete disintegration of the entire stone.\nIn addition to the compressive forces and the negative tensile forces, cavitation microjets contribute to calculus fragmentation. Acoustic cavitation occurs when the tensile forces exceed ambient pressure, pulling apart the liquid and creating a bubble that collapses with the return of positive pressure.\nB. Lithotripter Designs\nShock waves are created by converting energy into an acoustic form. The currently available extracorporeal shock wave lithotripters use energy sources that are electrohydraulic, piezoelectric, electromagnetic, or explosive in nature. Shock waves are focused on the stone by ellipsoidal reflectors, shaped array, or lens. Localization of the stone is by x-ray and/or ultrasound studies.\nWith the electrohydraulic method, electricity is discharged into water across a gap between two electrodes. The temperature of the water rises rapidly to form steam and then a plasma. A compressive pressure pulse results from expansion of the heated gases, followed by a negative pressure pulse as the gas bubble collapses. In the piezoelectric machines, an electrical field is applied across a piezoelectric crystal, changing the external dimensions of the crystal. Pressure waves are produced by the movement of the crystal. Multiple crystals are used in the machines for reliability and ease of construction. In the electromagnetic generators, a magnetic field is generated by current flow through a wire or coil. Magnetic materials are attracted or repelled by this field, turning electrical energy into mechanical and acoustic energy. The pressure wave is created by movement of a flexible membrane from passage of current through a fixed coil.\nFocusing is accomplished geometrically by an ellipsoidal reflector in many machines. The shock waves created at the first focus of the ellipsoid by the generation system are reflected by the ellipsoidal reflector to arrive at the second focus simultaneously creating a shock wave at the final focal point. In the Siemens electromagnetic lithotripter, the sound waves are focused by biconcave acoustic lens. In the piezoelectric lithotripters, the crystals are shaped as part of a sphere to focus the energy.\nIn order for the shock wave to effectively destroy the stone, the site of the stone must be coincident or nearly coincident with the final focal point of the generated shock wave. This is accomplished by first, localizing the site of the stone in relation to the geometric focal point of the shock wave and, then, targeting (or adjusting) the location of the stone so as to be coincident with the focal point of the shock wave.\nThe targeting of the stone in relation to the focal point of the shock wave is conventionally accomplished by physically adjusting the position of either the patient or the shock wave generating system, until the location of the stone is determined to be coincident with the geometric focal point of the shock wave generating system. Although somewhat cumbersome, this method of adjustment has proven useful in conventional lithotripsy.\nDifficulty has existed in locating the stone and accurately identify that location in relation to both the geometric focal point and, more importantly, the acoustic focal point of the machine. Stones are typically located using ultrasonography or x-ray flourography, with the image of the stone being displayed on a videoscope or x-ray image. A computer generated image, denoting the geometric focal point of the shock wave, may be superimposed on the ultrasound video display or x-ray image as a targeting aid. However, this technique of localizing and targeting has proven inadequate in several respects.\nUltrasound has the advantage of being able to be used continuously, but the disadvantage that stones cannot be visualized in certain locations. Moreover, ultrasound requires a considerable learning curve and sometimes it is difficult to determine the degree or fragmentation, in part due to the formation of an echogenic fragment line. Stone localization by ultrasound may be difficult in the kidney in the presence of a percutaneous nephrostomy tube, multiple calculi in the renal pelvis, or partial staghorn calculi. X-ray fluoroscopy is less commonly used and present added health concerns in that some people want strongly to avoid any exposure to x-rays. In addition, fluoroscopy is unable to visualize radiolucent stones such as those common in the gallbladder.\nMore importantly, the conventional techniques for localizing the stone frequently do not allow effective targeting. In order for the lithotripsy to effectively destroy the stone, the target center or the stone should be positioned as closely as possible to the focal point of the lithotripter. To the end, it is considered advantageous to be able to keep the stone in a region where the peak pressure at the stone site during lithotripsy is at least 50% of the pressure at the focus. For currently available lithotripters, the area in the focal plane described by this 50% isobar can be as small as 3 mm.sup.2, as reported by Coleman and Saunders (Ultrasound in Medicine and Biology, 15(3):213-227. (1989)). Conventional techniques of ultrasonography and flourography are problematic when targeting with such a small focal area for a variety of reasons, e.g., imprecision in manually adjusting the physical position of the patient and/or the shock wave focal point, displacement of the acoustic focal point of the generated shock waves from the geometric location due to diffraction effects, and stone visualization errors introduced due to diffraction effects on ultrasound imaging systems. As a result, extracorporeal lithotripsy procedures currently involve the application of between 1500 and 4000 shock waves to the stone site in order to reduce stone fragments to a size that can safely pass through the patient's system, or to be treated successfully using dissolution agents. Such a large number of shock waves prolong the treatment, may require multiple treatments, increase treatment cost, can produce tissue damage, and may require the administration of anesthesia to the patient.\nIn seeking to reduce the number of shocks applied, conventional targeting techniques have evolved along two diverging paths. One philosophy of improving lithotripters has been to apply more powerful shook waves and increase the volume of the acoustic focal point of the shock waves. This approach has the disadvantage of applying more shock wave energy to the tissue near the stone, with the attendant risk of tissue damage. The other philosophy of lithotripter development is to decrease the shock wave amplitude and decrease the volume of the focal point, the argument being that this method reduces the risk of tissue damage while maintaining the ability to destroy stones. A further advantage of this philosophy is that it makes possible anesthesia-free lithotripsy. However, a reduced focal point volume exacerbates existing targeting problems.\nOther improvements in localizing and targeting stones have been recently reported. See Kuwahara et al, \"Initial Experience Using a New Type Extracorporeal Lithotripter With An Anti-Misshot Control Device,\" J. of Lithotripsy and Stone Disease, Vol. 3, No. 2, pp. 141-146 (1991). This device is based noon the use of piezoelectric transducers to generate the shock wave, but the piezoelectric transducers are also used in a rudimentary way to assist in localizing the stone. Initial localizing of the stone is performed using an ultrasound sonoprobe installed in the shock wave generator. Ultrasound pulses are directed from the shock wave generator toward the geometric focal point of the shock wave. The sound levels, if any, reflected from the focal point region are detected by the piezoelectric transducers (now acting as a microphone). The signal is judged as a hit or miss based upon the amplitude of the reflected wave, i.e., a high level indicates a hit and a low level indicates a miss. A predetermined threshold is established in a comparator to judge whether the signal is a hit or a miss.\nSimply stated, the Kuwahara et al. technique merely provides \"yes/no\", trial-and-error, information on whether the stone has been accurately sited and targeted. Kuwahara et al. failed to recognize or disclose any guidance adjustment mechanism for determining the direction and amount of movement which is needed to accurately locate and target the stone. Moreover, the technique of Kuwahara et al. is limited to lithotripters which are based upon piezoelectric shock wave generators. Kuwahara et al. does not disclose a guidance adjustment mechanism which is useful in other types of lithotripters."}
{"text": "In the case, particularly, of a chronograph mechanism, it is known that exerting too great a force on a control member, such as a push-button, can act too violently upon the functions associated therewith. Thus, for example, it can happen that, on a crown or push-button used for operating a chronograph, the violence of the force transmitted to the chronograph mechanism damages it, which may mean that parts have to be changed.\nTo avoid such problems, it is known to protect push-buttons or crowns by using shoulders to border them or even using a fixed or removable frame, which completely or partially surrounds them. This prevents a shock from accidentally moving the push-buttons or crowns so as to protect the mechanism. However, the timepiece case then becomes complex and much more cumbersome. Incidentally, it may also become unattractive, which is contrary to the desired effect for mid range or top range timepieces."}
{"text": "All musical instruments, especially multi-stringed ones such as the piano and harpsichord, require periodic tuning to insure that they consistently reproduce the proper pitches when played.\nAt present this periodic tuning is accomplished primarily by listening for the beat signal between the note and an audible reference frequency. This is inherently an inaccurate method, and in the case of multi-stringed instruments such as pianos and harpsichords it requires skill far beyond that of the average musician.\nThere exist today two types of electronic aids which enable the person of no special skill to successfully tune a multi-stringed instrument. In the first of these aids the audible note from the instrument is converted into pulsations of light which are used to illuminate a rotating strobe disc. Deviation of the frequency of the note from the proper value is displayed as apparent movement of a pattern of light and dark spots on the disc.\nThe second method employs an electronic frequency counter to measure the frequency of the note and provide a digital display of the frequency in hertz. The user must then compare this frequency to the desired frequency and adjust the instrument accordingly.\nBoth these methods suffer from inherent disadvantages which have prevented their widespread use. The strobe disc method employs an electric motor and is therefore bulky as well as being relatively inaccurate. The frequency counter method has a slow response time, limiting its usefulness to the middle and upper octaves, in addition to which the digital readout is difficult for an untrained person to interpret.\nThe invention herein described seeks to overcome these disadvantages, thus making it possible for an untrained person to quickly tune any note in the musical spectrum with an accuracy exceeding that of the best professional ear."}
{"text": "Autocorrelation is a classic technique for finding signals in the presence of noise. Prior art autocorrelation techniques involve delaying the signal by predetermined increments, multiplying the delayed signal by the incoming signal, averaging the result, and repeating the above steps for predetermined time delays. The prior art requires thousands of multiplications for every time period.\nFor certain applications, for example, applications relating to lasers, time periods are short, on the order of 50-100 nanoseconds. In such applications, the number of multiplications required for each time period exceeds the capabilities of prior art systems. The motivation of the invention is to provide a simpler, more cost effective autocorrelation method and apparatus which works as well as the prior art."}
{"text": "1. Field of the Invention\nThe invention relates to the field of inspection devices and, in particular, to a self centering device that allows a laser measurement device to inspect the countersink in countersunk holes.\n2. Description of Related Art\nFlush head fasteners require that the fastener holes be countersunk. The inspection of countersink portion of holes is a time consuming operation; particularly if mechanical or optical inspection techniques are used. This is particularly true on large aircraft structures with a tremendous number of such fasteners. While automated drilling equipment greatly reduces the chance of out of tolerance fastener holes, inspection is still required. One of the most common methods is to use a Trulok gage, which consists of a spring loaded ball, which is seated in the countersink. It can determine only the depth of the countersink. It is calibrated assuming a 100 degree countersink. Several different balls are required because of the necessity of seating the balls tangent to the countersink. Finally it can not be used on complex surfaces. What is needed is a simple and automatic measuring device that can be installed in an automated fastener drilling machine or operated by hand.\nLasers are used in numerous inspection and measurement systems. One particularly useful device is a portable gap measurement system; particularly useful in inspecting the gaps between body panels and doors on automobiles. This device is a hand held and passed along the gap and automatically measures and records the gap and determines whether the gap is out of tolerance. The device consists of a laser transmitter, which passes a laser beam back and forth across the gap (forming a visual line across the gap) and an optical light receiver and recorder. The change in time of the received laser light as it passes over the edges of the gap is recorded by a laser line sensor which enables extremely high resolution profile scanning of up to 500 points across a horizontal measuring line of one inch length as well as the change when the laser light beam passes out of the gap. It is a simple calculation to determine the width of the gap and whether it is in tolerance or not. If such a device could be properly positioned over a countersunk hole it could be used determine if it was in tolerance or not.\nThus, it is a primary object of the invention to provide a device that can self center the laser beam to measure the dimensions of a countersink hole.\nIt is another primary object of the invention to provide a device that can self center the laser beam to measure the dimensions of a countersink that is portable.\nIt is a further object of the invention to provide a device that can measure the dimensions of a countersink that is portable based on existing gap measurement devices using laser beams."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor structure. More particularly, the present invention relates to a thin film transistor (TFT) having an oxide semiconductor material.\n2. Description of Related Art\nBy virtue of the rising awareness of environmental protection, flat panel displays featuring superior properties including low power consumption, optimal space utilization, no radiation, and high definition have gradually become the mainstream products in the market of the displays. Common flat panel displays include liquid crystal displays (LCDs), plasma displays, organic electroluminescent displays, and so forth. The prevailing LCD is mainly comprised of a thin film transistor (TFT) array substrate, a color filter substrate, and a liquid crystal layer sandwiched between the two substrates, for instance. Amorphous silicon (a-Si) TFTs or low temperature polysilicon TFTs on a conventional TFT array substrate often serve as switch elements of sub-pixels. According to the recent researches, the oxide semiconductor TFT, in comparison with the a-Si TFT, is characterized by favorable carrier mobility and, in comparison with the low temperature polysilicon TFT, is characterized by desired threshold voltage (Vth) uniformity. Hence, an oxide semiconductor TFT has the potential to be the key element in the next generation flat panel display.\nHowever, under the atmospheric environment, the oxide semiconductor material is apt to be affected by moisture and oxygen, such that electrical properties of the oxide semiconductor material are varied together with time. As such, electrical performance and reliability of the oxide semiconductor TFT are deteriorated."}
{"text": "1. Field of the Invention\nThe present invention relates to a circulation type hot water supply device that circulates hot water through a hot water supply heat exchanger of a water heater even during a stop of delivery of hot water, and has a function of diagnosing abnormalities of the heat exchanger at early stages.\n2. Description of the Related Art\nA conventional water heater comprises a hot water supply heat exchanger, a burner that heats water supplied to the heat exchanger from a water supply channel upstream of the heat exchanger, and delivered hot water temperature detection means for detecting a temperature of hot water fed from the heat exchanger to a hot water delivering channel downstream of the heat exchanger. A combustion amount of the burner is controlled so that a delivered hot water temperature detected by the delivered hot water temperature detection means reaches a set hot water temperature set by a remote controller.\nAnother conventional water heater comprises a flow rate sensor and a flow rate control valve provided in a water supply channel, and supplied water temperature detection means for detecting a temperature of water supplied through the water supply channel to a heat exchanger. In this water heater, a combustion amount required for increasing a delivered hot water temperature to a set hot water temperature is calculated from a deviation between a delivered hot water temperature detected by delivered hot water temperature detection means and a supplied water temperature detected by the supplied water temperature detection means, and a water supply amount detected by the flow rate sensor, and a combustion amount of a burner is controlled by feedforward control. If the delivered hot water temperature is not increased to the set hot water temperature even with a maximum combustion amount of the burner, the water supply amount is reduced by the flow rate control valve.\nImmediately after opening of a hot water delivering tap connected to a hot water delivering channel, cold water remaining in the hot water delivering channel between the heat exchanger and the hot water delivering tap flows out from the hot water delivering tap, which causes discomfort to a user. A conventional circulation type hot water supply device for solving such a problem is also known. In such a hot water supply device, a downstream end of a hot water delivering channel is connected to a water supply channel, and a circulating pump that returns hot water fed from a heat exchanger to a hot water delivering channel to the heat exchanger through the water supply channel is provided in the hot water delivering channel. In this device, water is supplied to the heat exchanger by an operation of the circulating pump even after the hot water delivering tap is closed to stop delivery of hot water, thereby causing combustion of a burner. This allows a temperature of the hot water in the hot water delivering channel to be maintained at an appropriate temperature, and allows the hot water at the appropriate temperature to be delivered immediately after the opening of the hot water delivering tap.\nIn areas with high water hardness, CaCO3 or MgCO3 contained in water is deposited in a heat absorbing pipe of a heat exchanger, and a deposit easily adheres to an inner surface of the heat absorbing pipe. Adhesion of the deposit reduces heat exchange efficiency of the heat exchanger to make it difficult for combustion heat of a burner to be absorbed by water, thereby increasing a temperature of the heat exchanger itself. Repeating combustion in this state causes heat damage to the heat exchanger and thus causes leakage.\nJapanese Patent Laid-Open No. 2003-254615 discloses a simple water heater that is not of a circulation type that can identify abnormalities of a heat exchanger at early stages based on an amount of increase of a detection temperature from delivered hot water temperature detection means after a hot water delivering tap is closed to stop combustion of a burner (so-called, an amount of delayed water temperature increase). If adhesion of a deposit to an inner surface of a heat absorbing pipe reduces heat exchange efficiency of the heat exchanger to increase a temperature of the heat exchanger itself as described above, the detection temperature from the delivered hot water temperature detection means is significantly increased after a stop of delivery of hot water. Thus, the presence of an abnormality can be determined when the amount of increase of the detection temperature from the delivered hot water temperature detection means after the stop of the delivery of hot water reaches a predetermined threshold value or more.\nIt is desired to identify abnormalities of a heat exchanger at early stages using such a technique in a circulation type hot water supply device. Some circulation type hot water supply devices always drive a circulating pump. In such a device, water is supplied to a heat exchanger by an operation of a circulating pump even after a hot water delivering tap is closed to stop combustion of a burner, thereby causing no delayed water temperature increase of the heat exchanger. Thus, even if a deposit adheres to an inner surface of a heat absorbing pipe, a detection temperature from delivered hot water temperature detection means is not increased after a stop of delivery of hot water, and abnormalities of the heat exchanger cannot be identified based on an amount of increase of the detected temperature.\nIn view of the above described circumstances, the present invention has an object to provide a circulation type hot water supply device that can reliably identify an abnormality of a heat exchanger when it occurs, such as adhesion of a deposit to an inner surface of a heat absorbing pipe."}
{"text": "This invention relates generally to shift registers and particularly to a shift register stage useful as a select line scanner for liquid crystal displays.\nLiquid crystal television and computer displays (LCDs) are known in the art. For example, see U.S. Pat. Nos. 4,742,346 and 4,766,430, both issued to G. G. Gillette et al. Displays of the type described in the Gillette patents include a matrix of liquid crystal cells which are arranged at the crossovers of data lines and select lines. The select lines are sequentially selected by a select line scanner to produce the horizontal lines of the display. The data lines apply the brightness (gray scale) signals to the columns of liquid crystal cells as the select lines are sequentially selected.\nPreferably, the drive circuitry, which drives the select line scanner, which selects the horizontal lines to be displayed, is fabricated directly onto the same substrate and at the same time as the liquid crystal cells are fabricated. Also, because a large number of data lines and select lines are required for a television or computer display, and because the small pixel pitch limits the space available for laying out the driver circuitry, it is essential to keep the circuitry as simple as possible.\nFIG. 1 illustrates an example of a known scan register described in U.S. Pat. No. 5,222,082, which may be integrated with a liquid crystal display device. This register is driven with multiphase clocking signals C1, C2, C3, with different ones of the clock phases applied to different ones of the scan register stages 11.\nFIG. 2 illustrates one of the scan register stages in detail. The scan register stage includes an input section including transistors 18 and 19, an intermediate section including transistors 20 and 21 and an output section including transistors 16 and 17.\nThe output section is arranged as a push-pull amplifier, with a clocked supply potential connected to its supply connection 14. An output is accessed at the interconnection of the transistors 16 and 17.\nThe input section is arranged as a switched amplifier to exhibit a predetermined potential during the clock phase applied to the supply terminal of the output section. The output signal, P1, of the input stage, is coupled to drive the output transistor 16. More particularly the output, P1, follows the input signal applied to the gate electrode of transistor 18. The output of the input section will be high when the clock phase applied to terminal 14 goes high, and a high level is translated to the output terminal 13. The high level at node P1 is retained at node P1 until the occurrence of a clock phase C3 when the input signal is low. Thus the gate of the output transistor 16 will be at a high level when the clock C1 goes high providing a charging path to output 13 and when clock C1 goes low, providing a path to discharge the output node 13.\nThe intermediate section is arranged as a clocked inverting amplifier responsive to the input signal. The output of the intermediate stage is coupled to the gate electrode of the pull down transistor 17 of the output stage. The intermediate stage includes pull up and pull down transistors 20 and 21 respectively. The conductance of transistor 21 is greater than that of transistor 20 so that if both transistors 20 and 21 are conducting concurrently, the output potential at node P2 will remain low. Thus if the clock applied to transistor 20 is high when the input signal is high, the output transistor 17 will be maintained in a non-conducting state. However since the application of the stage is as a scan register, input signal pulses occur relatively infrequently. As a result node P2 will normally be charged high for every clock pulse of clock phase C3 and output transistor 17 will normally be conducting.\nThe drains of transistors 18 and 20 receive a relatively positive biasing voltage V.sub.DD of about 16 volts. Thus node P2 is normally biased at about 16 volts. This places excessive stress on the gate electrodes of transistors 19 and 17 which tends to cause a considerable rise in their respective threshold voltages over time. As the threshold of transistor 19 increases, its ability to discharge node P1 decreases, and more time is required to turn off transistor 16. The result is that some of the clock C1 voltage may leak onto the output node 13 and undesirably affect subsequent register stages as well as erroneously addressing an LCD row of pixels.\nThe present invention addresses these problems and provides a shift register stage which not only precludes false output values, but also consumes less power."}
{"text": "Oil and gas extraction pipelines generally include joints between sections of pipe and between the main pipe and bypass pipes. The flow of oil or gas within such pipes can cause vibrations within the pipes and relative movement between pipes at joints. These vibrations and movement can ultimately lead to damage or failure of the joint. It is desirable to be able to monitor the physical conditions (such as strain) experienced at pipe joints. These measurements should ideally be available in real time in order that conditions causing high loads and potential failure can be identified. Current practice within the oil and gas industries, includes monitoring pipe joints by visual inspection using remotely operated vehicles (ROVs) or using intelligent pigs travelling within the pipes themselves. These pigs can measure structural parameters such as corrosion, cracks and pipe thickness. Both ROVs and pigs can only provide intermittent measurements, and inspection using these methods is costly. Pipe joints are also monitored using electrical strain gauges, inclinometers and accelerometers located on the pipes."}
{"text": "A computer network is a collection of interconnected computing devices that exchange data and share resources. In a packet-based network, such as the Internet, the computing devices communicate data by dividing the data into small blocks called packets, which are individually routed across the network from a source device to a destination device. The destination device extracts the data from the packets and assembles the data into its original form. Dividing the data into packets enables the source device to resend only those individual packets that may be lost during transmission.\nCertain devices within the network, such as routers, maintain routing information that describes routes through the network. Each route defines a path between two locations on the network. From the routing information, the routers may generate forwarding information, which is used by the routers to relay packet flows through the network and, more particularly to relay the packet flows to a next hop. In reference to forwarding a packet, the “next hop” from a network router typically refers to a neighboring device along a given route. Upon receiving an incoming packet, the router examines information within the packet to identify the destination for the packet. Based on the destination, the router forwards the packet in accordance with the forwarding information.\nLarge computer networks, such as the Internet, often include many routers that exchange routing information according to a defined routing protocol, such as the Border Gateway Protocol (BGP). When two routers initially connect, the routers exchange routing information and generate forwarding information from the exchanged routing information. Particularly, the two routers initiate a routing communication “session” via which they exchange routing information according to the defined routing protocol. The routers continue to communicate via the routing protocol to incrementally update the routing information and, in turn, update their forwarding information in accordance with changes to a topology of the network indicated in the updated routing information. For example, the routers may send update messages to advertise newly available routes or routes that are no longer available.\nIn the event one of the routers of a routing communication session detects a failure of the session, i.e., the session “goes down,” the surviving router may select one or more alternative routes through the network to avoid the failed router and continue forwarding packet flows. In particular, the surviving router may update internal routing information to reflect the failure, perform route resolution based on the updated routing information to select one or more alternative routes, update its forwarding information based on the selected routes, and send one or more update messages to inform peer routers of the routes that are no longer available. In turn, the receiving routers update their routing and forwarding information, and send update messages to their peers. This process continues and the update information propagates outward until it reaches all of the routers within the network. Routing information in large networks may take a long period of time to converge to a stable state after a network fault due to temporary oscillations, i.e., changes that occur within the routing information until it converges to reflect the current network topology. These oscillations within the routing information are often referred to as “flaps,” and can cause significant problems, including intermittent loss of network connectivity and increased packet loss and latency.\nTo reduce the impact of failures, some routers include a primary routing control unit and a secondary routing control unit. In the event the primary routing control unit fails, the secondary routing control unit assumes the responsibility of forwarding packet flows. During failover from the primary routing control unit to the secondary routing control unit, a significant period of time may elapse before the secondary routing control unit reaches a state in which it is able to process and forward packets. For example, the secondary routing control unit may need to reestablish routing communication sessions, e.g., BGP sessions, that were lost when the primary routing control unit failed. During this period, network traffic may be queued or lost.\nAs another technique for reducing the impact of failures, the failed router may also support “non-stop forwarding,” which refers to the ability to continue forwarding packets while the routing session is reestablished. Redundant components in the failed router maintain forwarding state information during control module failure, enabling the failed router to continue forwarding packets over routes that were available in the network's last-known state. Concurrently, the failed router relearns the network topology and recalculates its routing information and forwarding information. As a result, impact on current packet flows through the network is reduced.\nSome routers support “graceful restart,” which refers to the capability of preserving forwarding information while restarting a routing communication session, e.g., a BGP session. When establishing a routing communication session, a router that supports graceful restart advertises the capability to neighboring routers and specifies a restart time. The restart time is the estimated time that it will take for the router to reestablish the routing communication session after failure of the previous session and may be, for example, approximately 120 seconds. Upon failure of the routing communication session, the surviving router preserves forwarding information based on the expectation that the failed router will reestablish the routing communication session shortly. In other words, the surviving router will maintain the failed router within a forwarding path of the surviving router in the event of a failure of the routing communication session. Likewise, the failed router preserves forwarding information in a state that existed prior to the failure. Consequently, the surviving router does not need to find alternative routes unless the failed router does not reestablish the routing communication session within the advertised restart time. As a result, the routing instability caused by routing flaps within the network may be reduced."}
{"text": "An ultrasound diagnosis apparatus has been used for tests and diagnosis for various body tissues such as the heart, liver, kidney, and mammary gland. Recently, an ultrasonic contrast agent (hereinafter, “contrast agent”) of an intravenous administration type has been commercially introduced, and the contrast echo method has been used. With the contrast echo method, minute bubbles (hereinafter, “microbubbles”) or the like are intravenously administered as a contrast agent to enhance the blood-flow signals, realizing clear observation of the blood flowing state.\nFurthermore, research has been conducted on application of ultrasound to medical treatment. Through such research, it has been confirmed that, when emitting ultrasound waves onto the subject, pores are temporarily formed in cells by mechanical action of cavitation, and that they make genes and medicinal substances easy to permeate into the cells. It has also been confirmed that, when microbubbles are injected, pores are temporarily formed in cells by mechanical action of cavitation of the microbubbles even if the ultrasonic sound pressure is low, and that they enhance the permeation of genes and medicinal substances into the cells. The phenomenon that ultrasonic irradiation improves the permeation of genes or medical substances into the cells is sometimes called as “sonopolation”.\nThe problem to be solved by the present invention is to present an ultrasound apparatus that can enhance the penetration of genes and medicinal substances into a specific site by use of microbubbles, and a controlling method and a non-transitory computer readable medium for such an ultrasound apparatus."}
{"text": "Ultrasonic readings taken from fiber reinforced composites, such as plastics, can be related to the elastic modulus, and hence the strength of the material. Research results dating back to the 1960's have shown that stressing fiberglass reinforced polymers (FRP) may result in decreasing the modulus of the material—thus reducing its strength. In service, the stresses applied to FRP typically have the same effect.\nAn early patent application regarding use of ultrasound for testing materials was made in 1940 by Dr. Floyd Firestone at the University of Michigan. This early patent, and most subsequent work using ultrasound identified that invisible inhomogeneities within materials could be detected. For most of the 74 years since this patent application, the focus of ultrasonic testing has been on metals. In the case of ultrasound, a pressure pulse is applied to a metal material and inhomogeneities are detected when a feature blocks some of the path of the ultrasounds—features that are parallel to the path direction are generally not detected.\nUse of fiber reinforced composites for structural applications has been pursued since the 1930's, and has seen significant changes in the polymers and fibers available. With the growth of commercial aircraft starting in the 1960's, many investigations were conducted into use of ultrasound to detect flaws and defects in composites. Because many fiber reinforced composites are made in layers, interfaces between layers often interrupt the path of the pressure pulses and show as features or possible defects for most ultrasonic techniques. Ultrasound is generally considered to be the most common non-destructive technology used for composite materials.\nIn the early 1960's, use of ultrasonic testing (UT) was already showing reliable results for finding flaws in metallic structures. One of the desirable attributes of this technique is that reliable data could be generated if only one side of the material under investigation was accessible. This meant that in addition to finding flaws or defects, the same techniques could be used to produce thickness records of reasonable accuracy. At the same time, use of composite materials such as glass reinforced thermoset plastics was being explored for a number of structural and corrosion-resistant applications. Starting in the mid 1960's, researchers started to examine uses of ultrasound with these fiber-reinforced composite materials.\nIn one study, ultrasonic pulses were applied to composites and the responses were received using acousto-ultrasonic devices, thus mixing the principles of ultrasound with acoustic emission testing. This process is generally referred to as “acousto-ultrasonic” because the forces applied to the specimen are from ultrasonic pulses, whereas for acoustic emission, the forces applied to the composite are from mechanical loads, such as pressures and weights. In both cases, the responses are received in real time by acoustical equipment. This work showed correlation between the attenuation of the signal transmitted through the full thickness of a laminate—across its layers—and its tensile strength parallel to its layers. This technique is the subject of two American Society for Testing and Materials (ASTM) standards—ASTM E 1495 Standard Guide for Acousto-Ultrasonic Assessment of Composite, Laminates and Bonded Joints2 (ASTM E 1495) and ASTM E 1736 Standard Practice for Acousto-Ultrasonic Assessment of Filament Wound Pressure Vessels3 (ASTM E 1736).]\nA method to employ these techniques is described in ASTM E 1736. In this Standard Practice, it is recommended that initial readings be taken from the vessel to be monitored after calibration to a reference standard and before it is put into service. After the unit has been in service for some time, the results of the initial readings are then compared to readings taken after the unit has been in service. Changes that have occurred in the modulus of the composite from corrosion, decay or mechanical loads will appear as changes in the results of the scan. While there is a relationship between acousto-ultrasonic results and the presence of detectable defects such as voids or delaminations and porosity, it is not certain as yet whether or not these defects are the cause of strength changes. Furthermore, in order to make use of any correlation, it is generally required that reference standards be available for each feature and condition that requires detection.\nMany users of structural composites can report that the structural capacity of the composite has reduced while it has been in service. There have been numerous investigations into this phenomenon, including proposed models of the causes of these changes.\nEmbodiments of the system and method described herein are intended to address at least one of the drawbacks of conventional systems and methods."}
{"text": "Numerous attempts have been made to establish a noise reduction scheme, which provides good listening comfort and good speech intelligibility at the same time.\nIn prior art document EP 1326479 a method of reducing noise in a signal is described whereby an input signal is supplied to an amplification unit where the input signal is subject to an auxiliary noise reduction algorithm, to generate an auxiliary signal. The auxiliary signal is used to determine a control input for the amplification unit; and the amplification unit is then controlled with the control signal, to generate an output signal with reduced noise. According to this prior art the method comprises the further steps of: detecting the presence and absence of speech utterances; and in the absence of speech, determining a noise magnitude spectral estimate; and in the presence of speech comparing the magnitude spectrum of the audio signal to the noise magnitude spectral estimate; calculating an attenuation function from the magnitude spectrum of the audio signal and the noise magnitude spectral estimate; and modifying the input signal by the attenuation function, to generate an output signal with reduced noise. This prior art technique has the problem, that it requires a speech pause in which to determine the noise magnitude spectrum. In a typical party noise situation such a pause may not be detectable and the system then has little or no clue as to how to calculate the right attenuation function."}
{"text": "When electronic components are tested, testers are frequently used which are equipped with the electronic components to be tested. These testers generate a wave-shaped test signal which is applied to the electronic component to be tested in each case. In dependence on this input signal, the electronic component generates a dynamic digital response in the form of a data record which can be checked in various ways.\nOne possibility for evaluating this data record is represented by the evaluation in the frequency domain, for example by means of a Fourier transform. However, this evaluation uses a large amount of computing power and cannot, therefore, be executed during the run time of the test.\nA further possibility for evaluating the data record generated by the electronic component is represented by the evaluation in the time domain.\nThe computing method of sine wave fitting represents one evaluation method conceivable in this respect. The sine wave fitting method, too, requires a considerable amount of computing power. The data acquired cannot be evaluated during the run time, either. Instead, the entire data record must first be acquired and then supplied to a postprocessing unit. This is associated with a long measuring time and high measuring costs.\nAn alternative method for evaluation in the time domain is represented by the use of digital filters, the complexity of which, however, is relatively high.\nThe use of such digital filters necessitates high amounts of circuit area. The circuit area needed for digital filters is frequently larger than the circuit area for the electronic components to be tested.\nFor these and other reasons, there is a need for the present invention."}
{"text": "1. Field of the Invention\nThe invention relates in general to an electronic device, and more particularly to a liquid crystal display and a gate driver thereof.\n2. Description of the Related Art\nBecause a liquid crystal display (LCD) apparatus has the low power consumption, the low generated heat, the light weight and the non-radiative property, LCD apparatuses are used in various electronic products, and gradually replace the conventional cathode ray tube (CRT) display apparatuses. The liquid crystal molecule of the LCD apparatus possesses a property that the liquid crystal molecule cannot be always driven by the same voltage with the same polarity, or otherwise the property of the liquid crystal molecule is damaged after a long time and cannot rotate in response to the variation of the electric field even if the voltage is cancelled. Thus, the polarity of the voltage must be changed every other period of time to prevent the property of the liquid crystal molecule from being damaged.\nAt present, the frequently seen liquid crystal display panels may be classified into a normal pixel array and a flip pixel array. The polarity of the normal pixel array is switched by way of 1+2 line inversion, while the polarity of the flip pixel array is switched by way of column inversion. However, the flip pixel array has the smaller aperture ratio, while the polarity switching of the 1+2 line inversion has the higher power consumption. Thus, it is a very important subject to consider the high aperture ratio and the low power consumption."}
{"text": "Over the decades, wireless communication systems have become more and more technologically advanced, with performance increasing in terms of smaller size, operation at higher frequencies and the accompanying increase in bandwidth, lower power consumption for a given power output, and robustness, among other factors. The trend toward better communication systems puts ever-greater demands on the manufacturers of these systems.\nToday, the demands of satellite, military, and other cutting-edge digital communication systems are being met with microwave technology.\nMany of these systems use mixers to multiply signals and translate frequency. Mixers are used in both transmitter and receiver applications. Examples of microwave mixers that are built for this purpose are disclosed in Maas, S., Microwave Mixers, 2nd Edition, Artech House, 1993.\nMicrowave mixers may be categorized by the technology used for construction. For example, microwave integrated circuits (MICs) typically include discrete semiconductor components for microwave applications. Monolithic microwave integrated circuits (MMICs) often incorporate semiconductor devices directly on the circuit substrates, also for microwave applications. An alternative type of MMIC includes ceramic substrates with attached beamlead devices. In either case, copper or other appropriate metal is incorporated into the circuitry.\nAnother class of mixers utilizes Lumped Element Technology. Baluns comprising wire-wound transformers provide relatively broad bandwidths and small size, but have an upper frequency limitation. In addition, Lumped Element Technology is labor-intensive and therefore costly to produce.\nTypical MIC mixers are single-layered or double-sided and incorporate Schottky diodes. These mixers are usually passive devices, which do not require DC bias. Such circuits are suspended on metal frames or packaged in housings having pins, leads, or other connectors. MIC mixers perform well at high frequencies and over wide bandwidths. Generally, size increases as frequency decreases.\nThick film MMIC mixers, on the other hand, typically integrate passive Schottky diodes on ceramic substrates. The substrates themselves may form a surface-mount interface requiring no additional packaging for connecting to other electronic components. Thus, thick film MMIC mixers are generally small relative to MIC mixers. However, thick film MMIC mixers usually operate over narrow bandwidths relative to MIC mixers.\nThin film MMIC mixers typically incorporate diodes or field-effect transistors (FETs) directly on silicon or gallium arsenide substrates. Thin film MMIC mixers are smaller than MIC mixers, and are available in die form, but are commonly packaged as surface-mount components. Although such mixers are capable of operating at high frequencies, they usually also operate over narrow bandwidths relative to MIC mixers. Wide bandwidth operation is possible, but development cost is high, with associated design and foundry costs.\nIn sum, present technologies have several shortcomings that the present invention seeks to overcome. The bandwidth provided by MMIC technology is typically limited, and the development cost is high. Lumped Element Technology has an upper frequency limitation, and is labor-intensive to produce. MIC technology produces circuits that are physically larger, and utilizes metal frames or housings that further increase the size of the packaging."}
{"text": "The present invention relates to schema-driven object alignment on a canvas. In particular, it relates to the use of alignment guides to enable a user to align objects on a canvas.\nDigital canvasses may be used to display graphical representations of charts to users. Digital canvasses may include objects or shapes displayed at different locations on the canvas to provide the graphical representation to the user. The user is typically enabled to rearrange objects on the digital canvas to change the graphical representation to be more to their liking. It is known that alignment guides may be used to help align the objects on a digital canvas based on conceptual horizontal or vertical grid lines of the canvas."}
{"text": "1. Field of the Invention\nMy invention relates generally to an adjustable box wrench and particularly to a box wrench that positively grips the workpiece in response to torque applied to a cam lever separately articulated from the wrench body.\n2. Discussion of the Related Art\nThe box-end or box wrench is known in the art for turning polygonal nuts and bolts located in tight and inaccessible locations. The box wrench usually includes a thin-walled polygonally-fluted box structure on one or both ends. Because this thin-walled box structure is sized to fit snugly over the generally polygonal nut or bolt workpiece, the workpiece can be engaged for torquing even where there is only a few millimeters of clearance.\nAdjustable wrenches, including adjustable box-end wrenches, are also known in the art. For instance, in U.S. Pat. No. 4,325,275, David S. Colvin discloses an adjustable open-end and box-end wrench that uses a pair of spaced pins and skewed slots to allow adjustment of a member that cooperates with a box-end to re-size the box structure for a range of workpiece dimensions. Although Colvin's box-end wrench can be used to engage a workpiece in tight spaces, he neither considers nor suggests \"gripping\" means for positively engaging a deformed workpiece and instead relies on the matching polygonal geometry of box structure and workpiece to transfer torque therebetween. Similarly, in U.S. Pat. No. 3,858,465, Erik Lind discloses an adjustable wrench employing an axially displaceable and lockable external sleeve to adjust the geometry of the box structure but neither considers nor suggests means for positively gripping a deformed workpiece. In U.S. Pat. No. 3,363,490, K. Maichen discloses a double-ended simultaneously-adjustable wrench that operates by manually turning an eccentric pin to move two members cooperating to form surfaces for receiving and torquing polygonal workpieces. Maichen also neither considers nor suggests means for positively gripping a deformed workpiece. Swiss Patent No. 386 948 issued to August Samuel Aegerter and U. K. Patent 251,544 issued to Andrew Arbuckle both disclose adjustable box-end wrenches that rely on the manual rotation of a cam to urge a sliding adjusting rod into a position that fixes the effective engagement dimensions of a box structure. Neither Aegerter nor Arbuckle consider or suggest means for positively gripping a deformed workpiece.\nThe common problem of transferring effective torque to deformed nuts and bolts is well-known in the art. When a nut is \"frozen\" onto a threaded bolt, the torque applied in attempting to remove it may deform the polygonal geometry of the outer nut surface to such an extent that a common box-end wrench (even an adjustable one) no longer properly engages the deformed surface to transfer torque. Rounded comers merely slip within the engaging box structure when torque is applied, accomplishing nothing.\nGripping pliers and wrenches are known in the art for torquing deformed workpieces. For instance, in U.S. Pat. No. 3,611,843, Joachim E. Engel discloses an adjustable socket wrench that has a gripping handle and a pair of relatively movable jaws, one of which is secured to the gripping handle and the other of which is cammed into engagement with its neighbor. Engel's handle includes a coarse adjusting member coupled to an axially movable jaw to permit adjustment of the socket dimensions by turning a threaded member. Engel's wrench increases the gripping force in the jaws responsive to an increase in the pivotal rotation of the jaws relative to the gripping handle arising from torque applied to the handle but cannot grip deformed workpieces in tight spaces. Similarly, in U.S. Pat. No. 4,174,646, Simon Cotler discloses a universal tool with gripping action and replaceable jaws that has a body with an opening adapted to interchangeably receive cassette-type work elements such as a box-end wrench element or the like. Cotler uses a cam-locking lever integrated into the universal tool to apply a predetermined gripping force on the workpiece. Disadvantageously, his universal tool employs a bulky structure to receive the cassette-type working element and is not suited for gripping polygonal workpieces in tight spaces. Moreover, contrary to Engel's advantageous feature of increasing gripping force proportionately with applied torque, Cotler's cam-locking scheme provides an unvarying gripping force that may be insufficient to retain the workpiece at high torque levels. Finally, in U.S. Pat. No. 2,486,523, P. E. Deschenes discloses a similar cam-locking adjustable gripping wrench for use with bottle caps.\nThus, although gripping wrenches are known in the art, including cam-operated and cam-locking gripping pliers and wrenches, no gripping wrench suitable for use in the tight spaces serviced by box-end wrenches has been known in the art until now. A wrench adapted to gripping and torquing polygonal workpieces in tight spaces could satisfy a clearly-felt need in the art. The related unresolved problems and deficiencies are clearly felt in the art and are solved by my invention in the manner described below."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a metal/ceramic bonding article having a ceramic substrate and a metal plate bonded to the ceramic substrate, and a method for producing the same. More specifically, the invention relates to a metal/ceramic bonding article wherein a metal plate of a Cu—Ni or Cu—Mn—Ni alloy serving as a resistive element is bonded to a ceramic substrate and which is used for an electronic member for resistance, such as a shunt resistive element, and a method for producing the same.\n2. Description of the Prior Art\nIn conventional electronic members for resistance, such as shunt resistive elements, an alloy plate serving as a sheet resistor precisely worked by a press working or the like is bonded to a ceramic substrate via a brazing filler metal, such as an Ag brazing filler, by means of the brazing and soldering (see Japanese Patent Laid-Open No. 11-97203).\nHowever, in such an electronic member bonded via a brazing filler metal, a worked sheet resistor must be arranged on a ceramic substrate at a precise position to be bonded thereto. In particular, when a very small resistor having a size of 10 mm□, such as a hunt resistive element, is bonded to a ceramic substrate, the displacement of the resistor is easy to occur, and it is difficult to bond one or a plurality of resistors to one ceramic substrate, so that there is a problem in that the displacement failure of the resistor is not only caused, but productivity also considerably deteriorates.\nIn addition, since the alloy plate serving as the resistor has a low solder wettability, there is a problem in that there are some cases where bonding failure may occur when a semiconductor device or the like is mounted on the alloy plate."}
{"text": "The present invention relates in general to a knitting machine, in particular to a synchronizing device for the selection of knitting needles in flat or straight and circular knitting machines of the type having two pairs of needle beds arranged in tandem, knitting carriage circulating in one direction above the needle beds, each carriage supporting a device for selecting the knitting needles according to a desired pattern, the selecting device including a set of selection cams movable between a neutral and an active position, a set of solenoids for controlling the position of respective cams, means including a pulse generator for controlling actuation of selected solenoids, a pin board arranged in the path of movement of the selection cams and guiding a set of selection pins, and blade beds arranged on the needle beds and each including a set of Jacquard-type sinking blades assigned to respective knitting needles and being controlled by the selection pins.\nFrom prior art it is known how to select by means of selecting devices arranged on knitting carriages sinking blades for an auxiliary needle bed so as to control the individual knitting needles. In order to insure correct time points of the generation of control pulses by the selection devices, there are provided synchronizing devices consisting of pulse generators and synchronizing rails. The synchronizing rails have been integrally formed in the needle- or sinking blade beds or can be arranged on the knitting machine separately. For example, according to the German publication DE-OS No. 2919369 there is provided a pulse generator in the form of a disc cooperating with at least one, usually with two, pulse generators. The disc is arranged above the needle beds and is driven by a transmission chain synchronously with the carriages. From the German Pat. No. 2,114,013 known are also optical means for generating synchronizing pulses. In practice, such prior art solutions have been successful in flat knitting machines because they guarantee an error-free selection of the knitting needles, and consequently a high knitting quality.\nHowever, in straight and circular knitting machines operating with two beds of knitting needles arranged one after the other and using circulating carriages, the application of prior art solutions of selecting devices is not possible without further measures. Due to the two independent pairs of needle beds with corresponding different guiding tracks for the carriages, there result higher mechanical tolerances leading to higher timing tolerances for synchronization. Consequently, limitation of operating speed of the straight and circular knitting machines using circulating carriage would result. In addition, in the prior art solution using a chain power transmission from a drive to the knitting carriage, undesirable irregularities in the run of respective knitting carriages may occur especially upon braking the knitting machine a certain amount of kickback takes place. In order to reduce errors resulting from this kickback, it has been necessary to use always two spatially staggered pulse generators in order to recognize the direction change. In the case of straight and circular knitting machines it would be necessary to further increase the number of pulse generators and further component parts, which would impair the reliability and the efficiency of the knitting operation. So far, no measures have been devised how to improve the situation in circular knitting machines."}
{"text": "Optical discs must be provided with a way of generating synchronization or clock information and tracking information if data is to be properly recorded and read. Conventionally, tracking information is provided throughout each block of each segment of a data track of an optical disc and the read beam of the optical drive is continuously monitoring the track to provide tracking information on a continuous basis. Continuous generation of information often causes the tracking servo feedback signal to be noisy due to the varying data pattern coexisting with the tracking data, often resulting in tracking errors.\nAn approach to providing a relatively noise free tracking information servo feedback signal is to provide clock information and tracking information in only a small portion, the synchro-servo zone, of each sector or data block, with the informational data being provided in the remainder of each block. With this approach, clock information and tracking information are provided only periodically and not at the time informational data is being read or recorded, whereby the signals generated are relatively free of noise. The generation of the tracking information servo feedback signal utilized in this approach is the subject of this invention."}
{"text": "It is well known that fluorosilicones are used as lubricant oils and base fluids for greases. They also show good rheological properties and resistance to degradation and oxidation. However, this lubricant behavior, which is satisfactory at relatively low temperatures (e.g., up to 100.degree. C.) tends to progressively worsen at higher temperatures (e.g., from 100.degree. to 200.degree. C. and higher). The prior art has taught the use of various additives, such as antioxidants and/or antiwear agents, in an attempt to extend the practical application temperature range of fluorosilicones. But, e.g., Braun et al., in an article entitled \"Silicone Lubrication of Porous Bronze Bearings\" (Lub. Eng., 32, 176-182, 1975) found some antioxidants to be ineffective. Further, some fatty ester antiwear agents were found to give similar unsatisfactory results.\nKim et al. (U.S. Pat. No. 3,629,115) suggested the use of fluorinated phenylphosphine as both an antioxidant and an antiwear additive in fluorosilicones. However, the antiwear properties of these compositions are not satisfactory in practice.\nKobzova et al., in an article entitled \"Efficiency of Antioxidants in Phthalocyanine Greases\" (Khim. i Tekhn. Topliv i Masel, No. 10, pp 59-61, 1971) showed that conventional antioxidants, such as N-phenyl-.alpha.-naphthylamine, are nearly ineffective in fluorosilicone-based greases.\nTherefore, from an application standpoint, there is still a need for agents which, when added to fluorosilicones, result in effective lubricants for use in air at temperatures higher than 100.degree. C., particularly at 125.degree. C. to 200.degree. C., and higher."}
{"text": "Automatic, high-speed check-processing/imaging systems are well known in the art. Many difficulties and disadvantages of present systems stem from how well the imaging-radiation is tailored to the detectors. Since the checks can be mis-imaged by poorly matching imaging radiation to its associated electronic camera, the result is apt to be errors--e.g. check-images that include optical \"noise\" or other distracting information, or that \"miss\" meaningful data.\nWorkers recognize that electronically-stored data (\"electronic images\") can be processed much quicker, more reliably and less subject to error. But to do so, one must first capture an accurate image, or modified image, of the physical document and convert this into electronic computer (EDP) signals. The EDP image-signals can then be manipulated (e.g. be reproduced for visual review, be sorted and distributed, etc.) much more rapidly, easily and reliably than physical documents.\nCurrent systems contemplated for \"electronic image-lift\" teach using a video camera by which an operator views the front and back of the actual document as desired. Based on what he sees, the operator can electronically enter document-data into a computer system; e.g., such things as check-amount, account number and other data necessary for processing document transactions. Such physical viewing is labor-intensive, is subject to error (e.g. from operator fatigue) and is substantially slower than an automated image-capture arrangement.\nWorkers are beginning to think of using imaging technology as a way of improving document processing, as disclosed, for example, in U.S. Pat. Nos. 4,510,619; 4,205,780; 4,264,808; and 4,672,186. Generally, imaging involves optically scanning documents to produce electronic images that are processed electronically and stored on high capacity storage media (such as magnetic disc drives and/or optical memory) for later retrieval and display. It is apparent that document imaging can provide an opportunity to reduce document handling and movement, since the electronic images can be used in place of the actual documents.\nIt would be somewhat conventional to think of document processing with \"image capture\" using video cameras, with two light sources, one to illuminate each side of a document, plus various lenses to focus light onto the document. Successive document-images (\"image slices\") can be reflected from the document, front and back, into respective video cameras. These can convert the optical image into electronic signals; which can then be converted by appropriate circuitry into digital form. But the foregoing would have serious disadvantages; e.g. it would require two light sources and two camera systems--something expensive to provide and cumbersome to coordinate.\nThis invention addresses such disadvantages; e.g. teaching use of a single, high-intensity, well-cooled light source (cf. high-output xenon bulb, requiring substantially less power than a two-lamp system); and mounting the light source and associated optical components on a single base, and under a document transfer track, for ready access (e.g. for maintenance) and for better thermal isolation. Also, the taught system is modular (e.g. to plug-in to a relatively conventional sorter); it simplifies service and manufacture using interchangeable, easily-installed components. The system disclosed uses randomly-distributed fiber optics and electronic image conversion means in a high speed, automated \"image-lift\" system; e.g. one capable of accommodating the present advanced needs of financial institutions for document processing.\nAn object hereof is to address at least some of the foregoing problems and to provide at least some of the mentioned, and other, advantages."}
{"text": "Technical Field\nThe disclosure concerns a chamber or reactor for processing a workpiece such as a semiconductor wafer using microwave power.\nDescription of Related Art\nProcessing of a workpiece such as a semiconductor wafer can be carried out using a form of electromagnetic energy, such as RF power or microwave power, for example. The power may be employed, for example, to generate a plasma, for carrying out a plasma-based process such as plasma enhanced chemical vapor deposition (PECVD) or plasma enhanced reactive ion etching (PERIE). Some processes need extremely high plasma ion densities with extremely low plasma ion energies. This is true for processes such as deposition of diamond-like carbon (DLC) films, where the time required to deposit some type of DLC films can be on the order of hours, depending upon the desired thickness and upon the plasma ion density. A higher plasma density requires higher source power and generally translates to a shorter deposition time.\nA microwave source typically produces a very high plasma ion density while producing a plasma ion energy that is less than that of other sources (e.g., an inductively coupled RF plasma source or a capacitively coupled RF plasma source). For this reason, a microwave source would be ideal. However, a microwave source cannot meet the stringent uniformity required for distribution across the workpiece of deposition rate or etch rate. The minimum uniformity may correspond to a process rate variation across a 300 mm diameter workpiece of less than 1%. The microwave power is delivered into the chamber through a microwave antenna such as a waveguide having slots facing a dielectric window of the chamber. Microwaves propagate into the chamber through the slots. The antenna has a periodic power deposition pattern reflecting the wave pattern of the microwave emission and the slot layout, rendering the process rate distribution non-uniform. This prevents attainment of the desired process rate uniformity across the workpiece.\nA limitation on processing rate is the amount of microwave power that can be delivered to a process chamber without damaging or overheating the microwave window of the chamber. Currently, a microwave window, such as a quartz plate, can withstand only low microwave power levels at which DLC deposition processes can require hours to reach a desired DLC film thickness. The microwave window provides a vacuum boundary of the chamber and is consequently subject to significant mechanical stress, rendering it vulnerable to damage from overheating."}
{"text": "The present invention relates to packages for power semiconductor chips, and, more particularly, to hermetic power chip packages.\nPower semiconductor chips (hereinafter, simply \"power chips\") generate waste heat during operation typically in excess of about one watt. This heat must be removed in order to prevent destruction of the power chip. Power chips are accordingly assembled in a package or arrangement designed to facilitate the removal of heat from the power chip. Hermetic, or airtight, packages are particularly desirable for packaging power chips, since hermetic packages shield the power chips from contaminants and moisture that are known to degrade the operating performance of power chips.\nA typical hermetic power chip package that is presently available includes a relatively massive metallic baseplate on which the power chip is thermally mounted, and, which, in turn, is adapted to be thermally mounted upon a metallic heat sink. Two glass-to-metal hermetic seals are typically incorporated in the power chip package to permit electrical access to the power chip via current leads. These glass-to-metal seals are expensive to make, and the use of the metallic baseplate is also expensive. The power chip package additionally includes a housing that hermetically encloses the power chip and one or more of the foregoing glass-to-metal seals resulting in further expense of the package.\nThe foregoing describes a hermetic power chip package in an essentially completed form, that is, in a form ready for use in circuit applications. Initial electrical testing of power chips has heretofore been carried out by assembling the power chips in completed hermetic power chip packages to provide for removal of waste heat. Such testing is necessary to ascertain important device characteristics, such as, in a power Darlington transistor, the common emitter current gain, H.sub.FE, and the collector-to-emitter voltage at device saturation, V.sub.CE(SAT). If the power chip in a power chip package does not meet the required standards, the entire power chip package is discarded. As a consequence, the testing of power chips in presently available hermetic packages is costly. In order to make testing of power chips more economical, it would be desirable to provide a hermetic power chip package in a partially-completed or building block form, so as to reduce the expense of discarded packages."}
{"text": "The present invention relates to baseball training. More particularly, the invention relates to a method for training a baseball trainee in the art and science of batting a baseball.\nAt present the training available for properly and expertly batting a baseball is restricted to either manual instruction or mechanical devices which are constructed and designed to teach by rote memory or by trial and error. There does not exist in the prior art any device which allows the trainee to actually feel the way he or she should swing a baseball bat for maximum advantage. But it is virtually impossible to teach the trainee efficiently by either prior-art method. What is needed is a method which will simulate the ideal swing, and which will simultaneously communicate this simulation as a feeling to the trainee. The present invention provides such a machine."}
{"text": "Disclosed herein are various embodiments that relate to a light-emitting part for a vehicle, and particularly, to a light-emitting part for a vehicle that includes a curved plate-shaped light-emitting portion that uniformly emits light.\nHitherto, light-emitting parts for vehicles have been used. Some known light-emitting parts for a vehicle include a plurality of prism portions disposed in a planar shape on a light guiding plate in order to provide surface light-emission without using a plurality of dot-shaped light sources. Light incident from the dot-shaped light sources to the light guiding plate is reflected in a planar shape by the plurality of prism portions toward a plate-shaped transmissive body. Further, some examples of a method of uniformly emitting light from the planar plate-shaped transmissive body include a method of changing the height of the prism portion or the gap between the prism portions in response to the distance between the light source and the prism portion in consideration of the light attenuation effect, for example, as described in Japanese Patent Document No. 2006-222030 A (“the '030 Document”) and Japanese Patent Document No. 2008-15448 A (“the '448 Document”).\nJapanese Patent Document No. 11-321478 A (“the '478 Document”) discloses a light-emitting part using a decorative sheet. A decorative sheet is used to allow light from metal to be transmitted therethrough so that a metallic effect is visually obtained. Specifically, an adhesive layer, an anchor layer, a metal deposition layer exhibiting a metallic effect, an anchor layer, a colored layer, a hard coat layer, and a base film having a surface subjected to hairline pattern processing are laminated from a surface of a resin layer used as a base and formed as a lamination film serving as a transfer foil. Then, a sheet is subjected to in-mold molding to transfer a pattern to the surface of the resin layer as the base, and the base film is peeled off after the in-mold molding to finish a surface of a molded product with a metallic effect.\nFurther, Japanese Patent Document No. 2004-90897 A (“the '897 Document”) discloses another example of a light-emitting part for vehicles. In the light-emitting part disclosed in the '897 Document, a vehicle interior panel is formed by an optically transparent panel having a design portion laminated thereon and has a desired design at a vehicle interior side surface so that a light shielding region, shielding the light from the opposite side to the vehicle interior side, and a light transmitting region, transmitting the light from the opposite side to the vehicle interior side, are alternately arranged, and a light source is disposed inside the optically transparent panel. Then, when the light source is lit at night, the light is incident from the rear surface of the optically transparent panel. Subsequently, when incident light enters the light transmitting region and is transmitted through the front surface of the optically transparent panel, the light illuminates the dark vehicle interior. Thus, the light-emitting part serves as a decorative illumination unit.\nWhen the height or the pitch of the prism portion is changed in response to the distance from the light source in order to uniformly emit the light from the planar plate-shaped transmissive body, the uniformity of the light emitted from the light guiding plate is maintained even when the light is emitted from the transmissive body if the gap between the light guiding plate and the transmissive body is constant. However, when there is a difference in gap between the light guiding plate and the transmissive body due to the combination of the planar plate-shaped light guiding plate and the curved plate-shaped transmissive body, an attenuation degree of the light is changed due to a difference in gap therebetween. Accordingly, even when the light emitted from the light guiding plate is uniform, the uniformity of the light is not kept when the light is emitted through the transmissive body.\nFor this reason, in the techniques disclosed in the above-cited references, the light can be uniformly emitted from the transmissive body when the gap between the light guiding plate and the transmissive body is constant. However, when the gap between the light guiding plate and the transmissive body is not constant, the light cannot be uniformly emitted from the transmissive body.\nFurther, there is a tendency to increase the complexity of the shape of vehicle interior parts in order to improve the design thereof. Particularly, there are many members of vehicle interior parts that use a curve or a curved surface to achieve an improved design. Such a member having such a curve or curved surface may give a psychologically calming impression to a passenger of the vehicle. However, when the sheet is disposed on the curve and the curved surface by using the decorative sheet (the film) of the '478 Document, a problem arises in that a tension difference occurs in the case of a complex shape. Particularly, this problem noticeably occurs due to a tension difference between the flat surface portion and the curved surface portion. In this way, when the tension difference occurs, a crack may occur in the high-tension portion, that is, the curved surface portion. Thus, there is a possibility that cracks or visual imperfections may occur on the design. The light-emitting part for the vehicle is obtained by using the decorative sheet (the film) as described above so that the light is emitted from the light source to the decorative sheet (the film). However, in the decorative sheet (the film) of the '478 Document, the decorative appearance is simplified since the layers are uniformly laminated as multiple layers. For this reason, there has been a demand for a technique of improving durability while ensuring a better design of the decorative film.\nFurther, in the light-emitting part for the vehicle of the '897 Document, a decorative lamination structure is described in which a light shielding layer, including a transmissive portion transmitting the light transmitted through the optical transparent panel, and a light shielding portion, shielding the light, and a design layer, disposed at the vehicle interior side in relation to the light shielding layer and formed to match the surface shape of the optical transparent panel, are laminated on the vehicle interior side surface of the optical transparent panel. According to the light-emitting part with such a configuration, a design such as a metallic effect or a grain effect of the design layer is visually recognized at daytime, and the light emitted from the light source is visually recognized from the design layer through the transmissive portion at nighttime. Accordingly, the design layer is visually recognized while the light is transmitted therethrough at nighttime in a manner differently from the daytime. Here, in order to improve the transmissivity of the light emitted from the light source toward the vehicle interior side for the decorative illumination function, it is desirable to form the design layer as thin as possible. However, when the film thickness of the design layer is thin, a boundary line between the light shielding portion and the transmissive portion forming the light shielding layer is visually recognized at the vehicle interior side when the light is not emitted therefrom. Accordingly, a passenger feels uncomfortable when the vehicle interior is bright, and hence a good design is not obtained. For this reason, there has been a demand for a light-emitting part for a vehicle which improves the functionality as decorative illumination and the texture of a decorative illumination region to obtain a good design."}
{"text": "Corticotropin releasing factor (herein referred to as CRF), a 41 amino acid peptide, is the primary physiological regulator of proopiomelanocortin(POMC)-derived peptide secretion from the anterior pituitary gland [J. Rivier et al., Proc. Nat. Acad. Sci. (USA) 80:4851 (1983); W. Vale et al., Science 213:1394 (1981)]. In addition to its endocrine role at the pituitary gland, immunohistochemical localization of CRF has demonstrated that the hormone has a broad extrahypothalamic distribution in the central nervous system and produces a wide spectrum of autonomic, electrophysiological and behavioral effects consistent with a neurotransmitter or neuromodulator role in brain [W. Vale et al., Rec. Prog. Horm. Res. 39:245 (1983); G. F. Koob, Persp. Behav. Med. 2:39 (1985); E. B. De Souza et al., J. Neurosci. 5:3189 (1985)]. There is also evidence demonstrating that CRF may also play a significant role in integrating the response of the immune system to physiological, psychological, and immunological stressors [J. E. Blalock, Physiological Reviews 69:1 (1989); J. E. Morley, Life Sci. 41:527 (1987)].\nClinical data has demonstrated that CRF may have implications in psychiatric disorders and neurological diseases including depression, anxiety-related disorders and feeding disorders. A role for CRF has also been postulated in the etiology and pathophysiology of Alzheimer's disease, Parkinson's disease, Huntington's disease, progressive supranuclear palsy and amyotrophic lateral sclerosis as they relate to the dysfunction of CRF neurons in the central nervous system [for review see E. B. De Souza, Hosp. Practice 23:59 (1988)].\nIn affective disorder, or major depression, the concentration of CRF is significantly increased in the cerebral spinal fluid (CSF) of drug-free individuals [C. B. Nemeroff et al., Science 226:1342 (1984); C. M. Banki et al., Am. J. Psychiatry 144:873 (1987); R. D. France et al., Biol. Psychiatry 28:86 (1988); M. Arato et al., Biol Psychiatry 25:355 (1989)]. Furthermore, the density of CRF receptors is significantly decreased in the frontal cortex of suicide victims, consistent with a hypersecretion of CRF [C. B. Nemeroff et al., Arch. Gen. Psychiatry 45:577 (1988)]. In addition, there is a blunted adrenocorticotropin (ACTH) response to CRF (i.v. administered) observed in depressed patients [P. W. Gold et al., Am J. Psychiatry 141:619 (1984); F. Holsboer et al., Psychoneuroendocrinology 9:147 (1984); P. W. Gold et al., New Eng. J. Med. 314:1129 (1986)]. Preclinical studies in rats and non-human primates provide additional support for the hypothesis that hypersecretion of CRF may be involved in the symptoms seen in human depression [R. M. Sapolsky, Arch. Gen. Psychiatry 46:1047 (1989)]. There is preliminary evidence that tricyclic antidepressants can alter CRF levels and thus modulate the numbers of CRF receptors in brain [Grigoriadis et al., Neuropsychopharmacology 2:53 (1989)].\nThere has also been a role postulated for CRF in the etiology of anxiety-related disorders. CRF produces anxiogenic effects in animals and interactions between benzodiazepine/non-benzodiazepine anxiolytics and CRF have been demonstrated in a variety of behavioral anxiety models [D. R. Britton et al., Life Sci. 31:363 (1982); C. W. Berridge and A. J. Dunn Regul, Peptides 16:83 (1986)]. Preliminary studies using the putative CRF receptor antagonist .alpha.-helical ovine CRF (9-41) in a variety of behavioral paradigms demonstrate that the antagonist produces \"anxiolytic-like\" effects that are qualitatively similar to the benzodiazepines [C. W. Berridge and A. J. Dunn Horm. Behav. 21:393 (1987), Brain Research Reviews 15:71 (1990)]. Neurochemical, endocrine and receptor binding studies have all demonstrated interactions between CRF and benzodiazepine anxiolytics providing further evidence for the involvement of CRF in these disorders. Chlordiazepoxide attenuates the \"anxiogenic\" effects of CRF in both the conflict test [K. T. Britton et al., Psychopharmacology 86:170 (1985); K. T. Britton et al., Psychopharmacology 94:306 (1988)] and in the acoustic startle test [N. R. Swerdlow et al., Psychopharmacology 88:147 (1986)] in rats. The benzodiazepine receptor antagonist (Ro15-1788), which was without behavioral activity alone in the operant conflict test, reversed the effects of CRF in a dose-dependent manner while the benzodiazepine inverse agonist (FG7142) enhanced the actions of CRF [K. T. Britton et al., Psychopharmacology 94:306 (1988)].\nThe mechanisms and sites of action through which the standard anxiolytics and antidepressants produce their therapeutic effects remain to be elucidated. It has been hypothesized however, that they are involved in the suppression of the CRF hypersecretion that is observed in these disorders. Of particular interest is that preliminary studies examining the effects of a CRF receptor antagonist (.alpha.-helical CRF.sub.9-41) in a variety of behavioral paradigms have demonstrated that the CRF antagonist produces \"anxiolytic-like\" effects qualitatively similar to the benzodiazepines [for review see G. F. Koob and K. T. Britton, In: Corticotropin-Releasing Factor: Basic and Clinical Studies of a Neuropeptide, E. B. De Souza and C. B. Nemeroff eds., CRC Press p221 (1990)].\nIn order to study these specific cell-surface receptor proteins, compounds must be identified which can interact with the CRF receptor in a specific manner dictated by the pharmacological profile of the characterized receptor. Toward that end, there is evidence that the direct CRF antagonist compounds and compositions of this invention, that can attenuate the physiological responses to stress-related disorders, will have potential therapeutic utility for the treatment of depression and anxiety-related disorders. All of the aforementioned references are hereby incorporated by reference.\nPCT Application US94/1105 teaches 1N-alkyl-N-arylpyrimidines and derivatives thereof in the treatment of affective disorders, anxiety, depression, post-traumatic stress disorders, eating disorders, supranuclear palsey, irritable bowl syndrome, immune supression, Alzheimer's disease, gastrointestinal diseases, anorexia nervosa, drug and alcohol withdrawal symptoms, drug addiction, inflammatory disorders, or fertility problems.\nU.S. Pat. No. 5,062,882 teaches the synthesis of aryloxy- and arylthiotriazines useful as herbicides.\nU.S. Pat. Nos. 4,427,437 and 4,460,588 describe the synthesis of aryloxy- and arylthiopyrimidines useful for the killing of internal parasites, especially trematodes and nematodes, in warm blooded animals, and/or as herbicides for inhibiting the growth of severely damaging or killing plants.\nU. S. Pat. No. 5,281,707 teaches the synthesis and utility of water-soluble aryloxy triazines, useful for the thermal and photochemical stabilization of polyamide fiber materials.\nThe compounds and methods of the present invention provide the methodology for the production of specific high-affinity compounds capable of inhibiting the action of CRF at its receptor protein in the brain. These compounds would be useful in the treatment of a variety of neurodegenerative, neuropsychiatric and stress-related disorders. It is further asserted that this invention may provide compounds and pharmaceutical compositions suitable for use in such a method. Further advantages of this invention will be clear to one skilled in the art from the reading of the description that follows."}
{"text": "Minimally invasive medical techniques are intended to reduce the amount of extraneous tissue that is damaged during diagnostic or surgical procedures, thereby reducing patient recovery time, discomfort, and deleterious side effects. One effect of minimally invasive surgery, for example, is reduced post-operative hospital recovery times. Because the average hospital stay for a standard surgery is typically significantly longer than the average stay for an analogous minimally invasive surgery, increased use of minimally invasive techniques could save millions of dollars in hospital costs each year. While many of the surgeries performed each year in the United States could potentially be performed in a minimally invasive manner, only a portion of the current surgeries use these advantageous techniques due to limitations in minimally invasive surgical instruments and the additional surgical training involved in mastering them.\nMinimally invasive robotic surgical or telesurgical systems have been developed to increase a surgeon's dexterity and avoid some of the limitations on traditional minimally invasive techniques. In telesurgery, the surgeon uses some form of remote control (e.g., a servomechanism or the like) to manipulate surgical instrument movements, rather than directly holding and moving the instruments by hand. In telesurgery systems, the surgeon can be provided with an image of the surgical site at a surgical workstation. While viewing a two or three dimensional image of the surgical site on a display, the surgeon performs the surgical procedures on the patient by manipulating master control devices, which in turn control motion of the servo-mechanically operated instruments.\nThe servomechanism used for telesurgery will often accept input from two master controllers (one for each of the surgeon's hands) and may include two or more robotic arms on each of which a surgical instrument is mounted. Operative communication between master controllers and associated robotic arm and instrument assemblies is typically achieved through a control system. The control system typically includes at least one processor that relays input commands from the master controllers to the associated robotic arm and instrument assemblies and back from the instrument and arm assemblies to the associated master controllers in the case of, for example, force feedback or the like. One example of a robotic surgical system is the DA VINCI® system available from Intuitive Surgical, Inc. of Sunnyvale, Calif.\nA variety of structural arrangements can be used to support the surgical instrument at the surgical site during robotic surgery. The driven linkage or “slave” is often called a robotic surgical manipulator, and exemplary linkage arrangements for use as a robotic surgical manipulator during minimally invasive robotic surgery are described in U.S. Pat. Nos. 7,594,912; 6,758,843; 6,246,200; and 5,800,423; the full disclosures of which are incorporated herein by reference. These linkages often make use of a parallelogram arrangement to hold an instrument having a shaft. Such a manipulator structure can constrain movement of the instrument so that the instrument pivots about a remote center of manipulation positioned in space along the length of the rigid shaft. By aligning the remote center of manipulation with the incision point to the internal surgical site (for example, with a trocar or cannula at an abdominal wall during laparoscopic surgery), an end effector of the surgical instrument can be positioned safely by moving the proximal end of the shaft using the manipulator linkage without imposing potentially dangerous forces against the abdominal wall. Alternative manipulator structures are described, for example, in U.S. Pat. Nos. 7,763,015; 6,702,805; 6,676,669; 5,855,583; 5,808,665; 5,445,166; and 5,184,601; the full disclosures of which are incorporated herein by reference.\nA variety of structural arrangements can also be used to support and position the robotic surgical manipulator and the surgical instrument at the surgical site during robotic surgery. Supporting linkage mechanisms, sometimes referred to as set-up joints, or set-up joint arms, are often used to position and align each manipulator with the respective incision point in a patient's body. The supporting linkage mechanism facilitates the alignment of a surgical manipulator with a desired surgical incision point and targeted anatomy. Exemplary supporting linkage mechanisms are described in U.S. Pat. Nos. 6,246,200 and 6,788,018, the full disclosures of which are incorporated herein by reference.\nWhile the new telesurgical systems and devices have proven highly effective and advantageous, still further improvements are desirable. In general, improved minimally invasive robotic surgery systems are desirable. It would be particularly beneficial if these improved technologies enhanced the efficiency and ease of use of robotic surgical systems. For example, it would be particularly beneficial to increase maneuverability, improve space utilization in an operating room, provide a faster and easier set-up, inhibit collisions between robotic devices during use, and/or reduce the mechanical complexity and size of these new surgical systems."}
{"text": "1. Field of the Invention\nThe invention relates to web applications over the internet, and in particular, to a disruption-free web application execution mechanism.\n2. Description of the Related Art\nA web application is a program providing services through the hyper text transport protocol (HTTP), allowing a client to access services on a remote server over the internet. One field utilizing web applications is mobile communications. For mobile communications, temporarily disconnecting and reconnecting to the internet from a remote service frequently occurs, because a client is usually connected through variable environments where transmission signal quality may be uncertain. Thus, various offline execution methods are proposed to ensure that a remote service can be provided without disruption.\nConventional offline services are mostly implemented by specific modified programs provided by the remote server, wherein offline functions are defined should disconnection occurs, such as defined offline functions provided by Google Gears and Adobe AIR. The disadvantage of the approach is that since the offline functions require a specific modified program to be executed, flexibility for clients is decreased and clients must install the specific modified program for compatibility. Thus, a more compatible mechanism is desirable."}
{"text": "Electrostatic copying machines are known which comprise a rotatable cylinder which bears a surface layer of photoconductive material and which, as it rotates, brings successive parts of its surface first to a charging station, at which an electrostatic charge is imparted to the surface by a corona generator, then to an exposure station at which an image of a master to be duplicated is projected by an optical system to produce on the surface a charge pattern corresponding to the image areas of the master, then to a developing station at which a toner powder is applied to the image areas of the surface and finally to a transfer station at which the toner powder is transferred from the cylinder to a transfer sheet, which is thereafter subjected to treatment to fix to it the transferred toner powder. Such machines are expensive owing to the cost of the optical system and the corona generator and to the large amount of electrical power required for operation of the corona generator.\nU.S. Pat. No. 2576047 discloses an electrostatic duplicator which uses a rotating cylinder carrying a metal master bearing an image of high electrical resistivity. It uses corona generators both to charge the image and to assist in transferring the toner to a continuous roll of transfer material. Each of the numerous corona generators used needs 6KV to operate it and corona generators are known to be difficult to maintain in a condition in which they apply an even charge to a master. This machine has not been put into use commercially."}
{"text": "1. Field of the Invention\nThis invention relates to a top-emission type organic EL (electroluminescent) panel having an organic EL device formed by sequentially laying a reflecting electrode, organic compound layers and a transparent electrode on a substrate where a thin film transistor is formed and a protective layer covering the organic EL device. More particularly, the present invention relates to an organic EL panel designed to prevent the thin film transistor from being degraded by light.\n2. Description of the Related Art\nMassive efforts are being paid to develop organic EL devices. The organic EL device driving system is roughly classified into a passive matrix system and an active matrix system. However, in order to realize a large screen with high definition, an active matrix system with use of a thin film transistor is being developed.\nWhile amorphous silicon and poly silicon that are being employed to form thin film transistors have light susceptibility and hence there is a problem that some of the characteristics of thin film transistors can be changed and/or degraded when they are hit by light coming from organic EL devices and external light.\nJapanese Patent Application Laid-Open No. 2002-108250 describes a technique of forming a light-shielding layer that directly covers the thin film transistors on a substrate.\nThere is another problem about organic EL devices that light striking the organic EL devices is reflected by the reflecting electrodes and some of the undulations of the devices to fall the contrast ratio of the organic EL panel.\nJapanese Patent Application Laid-Open No. 2006-236963 discloses a technique of arranging a black adhering member on the device separation layer so as to expose the EL devices of an organic EL panel.\nSince Japanese Patent Application Laid-Open No. 2002-108250 describes a technique of forming a light-shielding layer that directly covers the thin film transistors, the technique can prevent some of the characteristics of thin film transistors from being changed or degraded by incident light.\nThe light-shielding layer may be formed by using a black inorganic insulating layer or a layer that is by turn formed by dispersing a black pigment or a black die in resin. However, the following problem can arise when such a light-shielding material is employed.\nSince the light-shielding layer is held in contact with thin film transistors, the light-shielding layer is sandwiched between the thin film transistors and an electro-conductive member when an electro-conductive member is arranged on the light-shielding layer. Then, if the light-shielding layer is formed by using a metal material or a material illustrating a large relative permittivity, the light-shielding layer is electrically charged to produce a capacitance, which by turn changes or degrades some of the characteristics of the thin film transistors.\nIn short, while the technique of Japanese Patent Application Laid-Open No. 2002-108250 can prevent some of the characteristics of thin film transistors from being changed or degraded by incident light, those characteristics of thin film transistors are changed or degraded due to the fact that a light-shielding layer is held in contact with thin film transistors in an organic EL panel.\nThus, there is not known any practical light-shielding technique for effectively covering thin film transistors.\nThe technique of Japanese Patent Application Laid-Open No. 2006-236963 can improve the contrast ratio of an organic EL panel by covering the region thereof other than the light emitting regions by means of a black material because the black material operates like a black matrix.\nHowever, since the device separation layer is tapered at the lateral surfaces thereof, which are then smoothly connected to the respective light emitting regions, external light that strikes the lateral surfaces of the device separation layer, which is part of light trying to enter the EL panel, is reflected and simply goes out. Therefore, when the lateral surfaces of the device separation layer are not covered by a black material but exposed, the contrast ratio of the organic EL panel can fall so much. In short, the technique of Japanese Patent Application Laid-Open No. 2006-236963 is not satisfactory for improving the contrast ratio of an organic EL panel because light striking the lateral surfaces of the device separation layer is reflected to become a cause of making the contrast ratio fall."}
{"text": "(a) Field of the Invention\nThis invention relates to the control of DC generators used in arc welding operations, the generators being driven by internal combustion engines.\n(B) Description of the Prior Art\nIn U.S. Pat. No. 2,135,046 issued on Nov. 1, 1938 is described a controller for a motor-generator arc welding set. This patent describes a motor powered by a three-phase electric main. The motor drives an A/C generator which supplies current to an arc welder circuit. A saturable transformer 65 has a primary winding 66 connected in series circuit relation with the welding cable and a secondary winding 64 connected in series with the coil of a relay which controls the main switch of the three phase main. The relay has a pair of contacts which are coupled to a timer 35 for the purpose of disconnecting the motor from the main after a predetermined interval of time following the cessation of welding operations.\nIn many arc welding machines, especially those used in the construction field wherein no electric energy is readily available, DC generators are used which are driven by internal combustion engines.\nThe problems associated with controlling a DC generator of a welding machine driven by an internal combustion engine are totally different from those encountered in connection with controlling an A/C generator driven by an A/C motor. In fact, the problems were compounded by the fact that the running or stopping of an internal combustion engine does not, in and by itself, produce an electric phenomenon which can be electrically monitored.\nAlso, the very high surge currents resulting from switching a DC welding generator into a short circuit condition discouraged workers in the art from developing a controller for DC arc welding machines, since the switching mechanism received in the welding circuit to implement the remote starting function becomes very easily damaged when such surge currents are allowed to flow during the switch closure. For these and other well-known reasons, no practical controller for remotely controlling the DC generator of an arc welder driven by an internal combustion engine has been introduced on the market, even though there is a great and pressing need for such a controller, especially since almost 40 years ago said U.S. Pat. No. 2,135,046 had described certain desirable features which such a controller might possess.\nThe present invention provides what is believed to be by applicants, the first practical, reliable and relatively inexpensive controller for arc welding machines whose DC generators are driven by internal combustion engines. This invention couples a vacuum-actuated device to the intake manifold of the engine for the purpose of carrying out switching operations. Also, applicants have discovered practical means for protecting the controller from the surge currents in the welding circuit."}
{"text": "1. Field of the Invention\nThe present invention relates to method and apparatus for communicating with one or more receiver stations disposed within a blackout region located in the vicinity of and surrounding each spot beam coverage area, the blackout region being caused by interference between concurrently transmitted and overlapping global coverage and spot coverage antenna radiated beams using the same frequency band. More particularly, this invention relates to method and apparatus for communicating with the one or more receiver stations disposed within the above-mentioned blackout region by selectively turning off the associated interfering spot beam during the predetermined interval of time to permit the global beam to communicate with the normally blacked out receivers.\n2. Description of the Prior Art\nIt is most desirable in a domestic satellite communications system to have a coexistence of overlapping spot and global coverage beams. For each high traffic ground station, for example, a separate spot coverage beam can be used for communication therebetween. Alternately, communication between the satellite and a plurality of low traffic ground stations can be accomplished with an overlapping global coverage beam where it might not be desirable to interconnect the individual low traffic ground stations to the nearest high traffic ground station for access to the satellite system. However, signals transmitted in the global beam cause interference in the spot beam signals and, therefore, to acquire an optimum signal in spot beam coverage area, various prior art interference suppression techniques may be employed.\nTo avoid signal degradation and permit separation of the overlapping spot coverage and global coverage beams, separate bandwidths or polarizations for the spot coverage beams and the global coverage beam can be used. Using separate bandwidths, however, results in inefficient use of the frequency spectrum and different polarizations may not be available where dual polarization is already used by each of the beams of the satellite system. In addition, substantial reduction of the interference between the overlapping global and spot coverage beams can be achieved, for example, through the use of antenna pattern shaping or channel coding techniques. More specifically, both antenna pattern shaping and channel coding techniques substantially remove interference in the spot beam coverage areas and dramatically reduce the area of interference in the vicinity of and surrounding these areas. However, neither prior art technique completely removes the interference in the area in the vicinity of and surrounding each of the spot beam coverage areas and, in fact, an area often termed \"a blackout region\" is produced within the above-mentioned area where signal separation between the two overlapping beams becomes virtually impossible without the use of highly sophisticated circuitry.\nNumerous other techniques have been devised to suppress interference particularly between two beams arriving at a receiver from separate directions. In this regard, see for instance, U.S. Pat. Nos. 2,579,882 issued to L. E. Thompson on Dec. 25, 1951; 3,094,695 issued to D. M. Jahn on June 18, 1963; 3,205,443 issued to D. R. Ludwig on Sept. 7, 1965; and 3,987,444 issued to R. J. Masak et al. on Oct. 19, 1976. Since the global and spot coverage beams transmitted from a satellite arrive at the area in the vicinity of and surrounding each spot beam coverage area from the same direction, techniques for separating signals arriving from different directions are not usable.\nThe problem still remaining in the prior art is to provide a technique which permits communication with receivers disposed within the blackout region formed in the vicinity of and surrounding each spot beam coverage area by the interference between signals in the concurrently transmitted and overlapping spot coverage and global coverage beams which use the same frequency band."}
{"text": "1. Field of the Invention\nThe present invention relates to a document illuminating system for use in an image reader such as a digital copying machine or an image scanner to illuminate an object such as a document as well as to an image reader incorporating such a document illuminating system. It is applicable to a digital copying machine, a facsimile machine, a printer, and a complex machine of these machines.\n2. Description of the Related Art\nIn recent years, LEDs (Light Emitting Diodes) have been well developed and highly bright LEDs are now available, LEDs offer such advantages as high longevity, high efficiency, high resistance to impact, and a mono-color light emission and will be expectedly adapted to various types of illumination including an illumination system of an image reader as a digital copying machine or an image scanner. Further, emission spectrum of a white LED covers a visible range of wavelength bands so that it is adoptable for a document illuminating system of a color image reader.\nFor this reason, various kinds of document illuminating systems using LEDs have been proposed. There is a well known technique to properly converge light beams with diffusion property from LEDs by a light guide plate or a reflective mirror in order to efficiently illuminate a document read area.\nFor example, Japanese Laid-open Patent Publication No. 2006-42016 (Reference 1) and No. 2006-295810 (Reference 2) disclose a document illumination system and an image reader to improve illumination performance of an illumination system by use of a reflective mirror. In both of the References, the position of the illumination system is diagonal relative to a document, the illumination system is configured to directly illuminate a document with a light component of a light beam from the LEDs at a small divergent angle to a direction orthogonal to a light emission face and to guide a light component at a large divergent angle to the document by use of a reflective mirror. Reference 1 exemplifies a structure with a reflective mirror including a curved face while Reference 2 exemplifies a structure with a reflective mirror including a planar face.\nHowever, there is a problem in the References 1 and 2 that the diagonal position of the illumination system to the document leads to an increase in thickness and size of the illumination system in an orthogonal direction to the document. In order to decrease the size of the reader, a distance from a light source to the document need to be decreased, which makes it impossible to illuminate the document with an even luminance distribution. Furthermore, this illumination system is configured to illuminate the document from a single direction so that dark lines may occur in an image in reading an uneven document.\nIn addition, Reference 2 discloses an LED structure having a special covering element to cover an emission face, however, such a covering element faces a problem of high manufacture cost since general low-cost LEDs are not suitable for such a structure.\nIn order to deal with occurrence of dark lines in an image, Japanese Laid-open Patent Publication No. 2008-67276 (Reference 3) discloses a structure to illuminate an object with light from two illumination elements which oppose to each other (hereinafter, sometimes referred to as opposite illumination). This can prevent dark lines in an image of an uneven document illuminated with light from a single direction.\nHowever, the two opposite illumination elements are each provided with light sources, which also results in increases in the size and manufacture costs of the illuminating system. From the sub-scan direction cross sectional (orthogonal to the document and including sub scan direction) drawing of the Reference 3, it is seen that the circuit board of light emitting elements extends diagonally downwards. Therefore, the document illuminating system of Reference 3 is increased in size in a height direction thereof. Moreover, since a normal line of the light emitting element directs to a direction of the document read area, it is difficult to secure a sufficient optical path length, leading to enlargement and uneven luminance of the illuminating system.\nIn view of solving the above problems, there is a demand for development of a document illuminating system in small thickness which can oppositely illuminates documents without occurrence of dark lines in an image.\nMoreover, along with an increasing demand for a higher-speed, higher image quality image reader, it is required to develop a document illuminating system configured to guide light onto the light-receiving face of the image sensor in even luminance distribution in the sub scan direction for the following reasons.\nA read line width of a CCD on a document corresponds to a size of the document in a main scan direction, and a position thereof in the sub scan direction may be shifted from a referential sub scan position due to an error in adjustment of such an optical element as a reflective mirror and an imaging lens or a movement of respective elements in reading a document from one end to the other in the main scan direction. Generally, the shift of a read line width may be about 0.2 to 1.0 mm in the sub scan direction relative to the optical axis of a read system.\nBecause of this, without illuminating a document with light amounts in even distribution, there will be an area insufficiently illuminated or an area over-illuminated in the document and a defective image (decrease in read accuracy) will occur. In order to prevent uneven illumination, a document area of about several mm in the sub scan direction need be illuminated in proper light amount.\nParticularly, a light receiving portion (document read area) of a digital copying machine or an image scanner is very narrow in the sub scan direction, about 0.1 mm, for example. Because of this, the document read area need be illuminated in proper amount such that the center of a luminance distribution curve comes at the document read area; otherwise, luminance of the area is substantially decreased.\nThere is a demand for a document illuminating system for use in a digital copying machine or an image scanner to illuminate the document read area with even luminance even when the center of illumination shifts from the document read area, so that a wide area of the image sensor in the sub scan direction receives light in proper amount with even luminance distribution.\nPreferably, the luminance distribution is to include, near a maximal value, an even portion in a width required to read a document with a fluctuation due to a mechanical error or the like added.\nThe even portion refers to a portion of the luminance distribution in which a data change rate of electric signals can fall within a practically allowable range by correcting the level of electric signals photo-electrically converted from an image by use of a signal processing circuit provided after the image sensor. For generating monochrome images, an input data change rate of about 30% is allowable by electrically correcting (amplifying) image signal values (gain adjustment). However, for generating color images, the input data change rate has to be about 12% or less since correction of color unbalance of three ROB colors need be taken into consideration rather than the correctable range of gain adjustment.\nHowever, References 1 to 3 cannot achieve even luminance distribution and need to be improved."}
{"text": "1. Field of the Invention\nThe present invention relates to IC testing devices for performing characteristic tests on ICs (Integrated Circuits) such as general-purpose AD converters, DA converters, and AD/DA converters with internal UCPUs microcomputers.\n2. Background Art\nFIG. 2 is a block diagram showing an example of the structure of a conventional IC testing device.\nIn the drawing, the analysis sample 1 is an IC which is the object of analysis, which in this example is a general-purpose AD converter.\nAdditionally, this testing device can be largely divided into an IC testing device main body 102 and two test stations (test station 103 and test station 104).\nIn this case, the IC testing device main body 102 is composed of a DC section 5, an LS section 6 and a solid-state relay 111. The DC section 5 measures the direct current characteristics of the analysis sample 1 based on analysis data received from test station 103 and test station 104. The LS section 6 selects the voltage level to be applied to the I/O terminals of analysis sample 1 for each terminal of the analysis sample 1.\nAdditionally, test station 103 and test station 104 are respectively composed of an analog circuit 7, a pin electronics section 8 and a power source 9.\nWhile the internal components of test station 103 and test station 104 are identical in the drawing, they are not necessarily identical.\nThe analog circuit 7 has a digital/analog converter (DAC), which DA converts analysis data (digital values) received from the IC testing device 102, and inputs the DA converted analog voltages into the analog input terminals (A.sub.in terminals) of the analysis sample 1 (AD converter).\nThe pin electronics section 8 compares digital values (digital values AD converted by analysis sample 1) outputted from the I/O terminals of the analysis sample 1 (AD converter) with reference values (digital values) received from the IC testing device main body 102.\nAdditionally, a tester bus which extends from the DC section 5 of the IC testing device main body 102 is serially connected in the order of analog circuit 7 of test station 103, pin electronics section 8 of test station 103, analog circuit 7 of test station 104 and pin electronics section 8 of test station 104.\nSimilarly, a tester bus which extends from LS section 6 of IC testing device main body 102 is serially connected in the order of the analog circuit 7 of test station 103, the pin electronics section 8 of test station 103, the analog circuit 7 of test station 104 and the pin electronics section 8 of test station 104.\nAdditionally, a solid-state relay 111 of the IC testing device main body 102 performs ON/OFF control of a 200 V AC supply to a power source 109 of each test station.\nSince the above-described conventional IC testing device shares the solid-state relay for performing ON/OFF control of the power source between the test stations, it is not capable of activating only one of the test stations.\nFor this reason, even in cases where a single test station is sufficient to perform the analyses due to reasons such as the number of ICs to be tested being few, both test stations must always be activated, thus wastefully expending electrical power.\nAdditionally, since the above-described conventional IC testing device serially connects the units of each test station by means of a tester bus, the flow of signals to the other test station is stopped in cases such as when one test station has failed, thus lowering the rate of operation of the device overall."}
{"text": "Embodiments of the present specification relate generally to wireless local area networks, and more particularly to systems and methods for real-time monitoring and characterization of transient network conditions experienced by client devices that are coupled to the wireless local area networks.\nA wireless local area network (WLAN) is a wireless computer network that links two or more components or devices using a wireless distribution method within a determined area such as a home, school, hospital, computer laboratory, or office building. Currently, WLAN hardware and protocols predominantly follow various versions of the IEEE 802.11 standard. In general, all components that can connect into a wireless medium in a WLAN are referred to as stations (STA). Typically, the stations are equipped with WLAN hardware, firmware, and drivers, also known as Wireless Network Interface Controllers (WNIC). These stations may be generally categorized as wireless access points (APs) and client devices. In one example, the wireless APs may include wireless routers which act as base stations for the WLAN. Some non-limiting examples of mobile and fixed client devices include patient monitoring devices, ambulatory hubs, wearable devices, mobile phones, tablets, laptops, handheld devices, and servers.\nWireless local area networks based on the IEEE 802.11 standard operate in several different modes. In one example, the client devices connect to the APs to join the WLAN. This enables users of the client devices to move around within a local coverage area and still be connected to the network. Moreover, the WLAN also provides a connection for the client devices to the wider Internet.\nTypically, the APs are strategically located at different locations within the local coverage area to provide the desired wireless service to client devices within range. Further, when a client device enters the WLAN, the client device may associate itself with one of the APs to communicate application data to an end device or another client device. Also, the client device may roam freely by transitioning from one access point (AP) to another AP within the coverage area of the WLAN.\nGenerally, in caregiving facilities such as hospitals, the client devices are primarily wireless telemetry devices. Wireless telemetry devices have a sensing component and a communication component that enable the measurement of physical or biological parameters and transmission of the measured data over a communication network. The performance and dependability of patient monitoring and telemetry is heavily dependent on the wireless local area network. Wireless telemetry devices are predominantly mobile, transitioning between APs within the coverage area of the hospital WLAN. Typically, device manufacturers of these wireless telemetry devices enter a service level agreement (SLA) with the hospitals to stipulate a minimal level of service that the wireless network must provide to a wireless telemetry device. However, enforcement of the SLA may be problematic because of the transient nature of the wireless links between client devices and various APs of the WLANs.\nCommonly, network performance with respect to the SLA may be determined at the AP by collecting network wide statistics. However, these statistics may only serve as a representation of general network health. Also, these statistics may fail to capture transient network characteristics in real-time at different coverage areas of the WLAN. In one example, statistics collected at the APs of the WLAN may capture the quality-of-service characteristics of the network. However, transient congestion points at certain network coverage areas may not be detectable at the APs. Moreover, these statistics may be limited to downlink traffic, overlooking key network metrics that may be only determined from client devices transmitting on the uplinks to the APs. Furthermore, leveraging the distributed location and mobility of the client devices within the wireless local network coverage area to discover transient network characteristics in real-time is often overlooked."}
{"text": "Rapamycin is a macrolide antibiotic produced by Streptomyces hygroscopicus which was discovered first for its properties as an antifungal agent. It adversely affects the growth of fungi such as Candida albicans and Microsporum gypseum. Rapamycin, its preparation and its antibiotic activity were described in U.S. Pat. No. 3,929,992, issued Dec. 30, 1975 to Surendra Sehgal et al. In 1977 Martel, R. R. et al. reported on immunosuppressive properties of rapamycin against experimental allergic encephalitis and adjuvant arthritis in the Canadian Journal of Physiological Pharmacology, 55, 48-51 (1977). In 1989, Calne, R. Y. et al. in Lancet, 1989, no. 2, p. 227 and Morris, R. E. and Meiser, B. M. in Medicinal Science Research, 1989, No. 17, P. 609-10, separately reported on the effectiveness of rapamycin in inhibiting rejection in vivo in allograft transplantation. Numerous articles have followed describing the immunosuppressive and rejection inhibiting properties of rapamycin, and clinical investigations have begun for the use of rapamycin in inhibiting rejection in transplantation in man.\nRapamycin alone (U.S. Pat. No. 4,885,171) or in combination with picibanil (U.S. Pat. No. 4,401,653) has been shown to have antitumor activity. R. Martel et al. [Can. J. Physiol. Pharmacol. 55, 48 (1977)] disclosed that rapamycin is effective in the experimental allergic encephalomyelitis model, a model for multiple sclerosis; in the adjuvant arthritis model, a model for rheumatoid arthritis; and effectively inhibited the formation of IgE-like antibodies.\nThe immunosuppressive effects of rapamycin have been disclosed in FASEB 3, 3411 (1989). Cyclosporin A and FK-506, other macrocyclic molecules, also have been shown to be effective as immunosuppressive agents, therefore useful in preventing transplant rejection [FASEB 3, 3411 (1989); FASEB 3, 5256 (1989); R. Y. Calne et al., and Lancet 1183 (1978).\nRapamycin has been shown to inhibit transplantation rejection in mammals [U.S. Pat. No. 5,100,899]. Rapamycin, its derivatives and prodrugs have also been shown to be useful in treating pulmonary inflammation [U.S. Pat. No. 5,080,899], systemic lupus erythematosis [U.S. Pat. No. 5,078,899], immunoinflammatory skin disorders, such as psoriasis [U.S. Pat. No. 5,286,730], immunoinflammatory bowel disorders [U.S. Pat. No. 5,286,731], ocular inflammation [U.S. Pat. No. 5,387,589], hyperproliferative vascular disorders, such as restenosis [U.S. Pat. Nos. 5,512,781 and 5,288,711], carcinomas [U.S. Pat. No. 5,206,018 and 4,885,171], and cardiac inflammatory disease [U.S. Pat. No. 5,496,832]; and in preventing the onset of insulin dependent diabetes mellitus [U.S. Pat. No. 5,321,009]. Additionally, rapamycin has been shown to be useful in treating adult T-cell leukemia/lymphoma [European Patent Application 525,960 A1], and in treating ocular inflammation [U.S. Pat. No. 5,387,589].\nBecause of its poor oil and water solubility, only a few formulations of rapamycin have proven satisfactory. U.S. Pat. Nos. 5,516,770 and 5,530,006 disclose intravenous rapamycin formulations, and U.S. Pat. Nos. 5,536,729 and 5,559,121 disclose liquid oral rapamycin formulations.\nMono- and diacylated derivatives of rapamycin (esterified at the 28 and 43 positions) have been shown to be useful as antifungal agents (U.S. Pat. No. 4,316,885) and used to make water soluble prodrugs of rapamycin (U.S. Pat. No. 4,650,803). Recently, the numbering convention for rapamycin has been changed; therefore according to Chemical Abstracts nomenclature, the esters described above would be at the 31- and 42-positions. U.S. Pat. No. 5,118,678 discloses carbamates of rapamycin that are useful as immunosuppressive, anti-inflammatory, antifungal, and antitumor agents. U.S. Pat. No. 5,100,883 discloses fluorinated esters of rapamycin. U.S. Pat. 5,118,677 discloses amide esters of rapamycin. U.S. Pat. No. 5,130,307 discloses aminoesters of rapamycin. U.S. Pat. No. 5,117,203 discloses sulfonates and sulfamates of rapamycin. U.S. Pat. No. 5,194,447 discloses sulfonylcarbamates of rapamycin.\nThe primary immunosuppressive agent presently used for inhibiting rejection in the allograft transplantation of organs in man is SANDIMMUNE (cyclosporine). Cyclosporine is a cyclic polypeptide consisting of 11 amino acids. The intravenous injectable formulation of SANDIMMUNE (IV) is a sterile ampoule containing, per ml, 50 mg of cyclosporine, 650 mg of Cremophor.RTM. EL and alcohol Ph Helv. (32.9% by volume) (under nitrogen). For administration this mixture is diluted further with 0.9% Sodium Chloride Injection or 5% Dextrose Injection before use. (Physicians' Desk Reference, 45th ed., 1991, pp. 1962-64, Medical Economics Company, Inc.) The macrolide molecule designated FK506, which has certain structural similarities to rapamycin, is also currently undergoing clinical investigation for inhibiting rejection in allograft organ transplantation in man. FK506 is isolated from Streptomyces tsuskubaensis and is described in U.S. Pat. No. No. 4,894,366 to Okuhara et al., issued Jan. 16, 1990 R. Venkataramanan et al., in Transplantation Proceedings, 22, No. 1, Suppl., 1 pp 52-56 (February 1990), report that the intravenous injectable formulation of FK506 is provided as a 10 mg/ml solution of FK506 in polyoxyethylated castor oil (HCO-60, a surfactant) and alcohol. The intravenous preparation must be diluted with saline or dextrose and administered as an infusion for 1 to 2 hours.\nThe Physicians' Desk Reference (45th ed., 1991, p. 2119, Medical Economics Company, Inc.) lists SANDIMMUNE (cyclosporine) as available in 25 mg and 100 mg strength capsules and as an oral solution in 50 ml bottles. The 25 mg capsules contain 25 mg cyclosporine, USP, and alcohol, USP dehydrated, at a maximum of 12.7% by volume. The 100 mg capsules contain cyclosporine, USP, 100 mg and alcohol, USP dehydrated, at a maximum 12.7% by volume. Inactive ingredients in the oral capsules are corn oil, gelatin, glycerol, Labrafil M 2125 CS (polyoxyethylated glycolysed glycerides), red iron oxide, sorbitol, titanium dioxide, and other ingredients. The oral solution is available in 50 mg bottles containing cyclosporine, USP, 100 mg and Ph. Helv. alcohol at 12.5% by volume dissolved in olive oil, Ph. Helv./Labrafil M 1944 CS (polyoxyethylated oleic glycerides) vehicle which must be diluted further with milk, chocolate milk or orange juice before oral administration.\nIMURAN (azathioprine, available from Burroughs Wellcome Co., Research Triangle Park, N.C.) is another orally administered immunosuppressive agent prescribed alone or in conjunction with other immunosuppressive agents. The Physicians' Desk Reference (45th ed., 1991, pp. 785-787, Medical Economics Company, Inc.) lists azathioprine as 6-[1-methyl-4-nitroimidazol-5-yl)thio]purine, which is provided for oral administration in scored tablets containing 50 mg azathioprine and the inactive ingredients lactose, magnesium stearate, potato starch, povidone, and stearic acid."}
{"text": "A satellite navigation system allows a user to determine his/her geographical position in three dimensions (longitude, latitude and altitude), his/her speed and the time. Satellite navigation systems are known under the names of GPS, GALILEO or GLONASS.\nSuch a satellite navigation system comprises a constellation of satellites in orbit. Each satellite emits a satellite signal consisting of a carrier with a determined frequency, modulated by a navigation message and an identification code called pseudo-random code. The navigation message notably contains data on the ephemeris of the satellite. The pseudo-random code is a pseudo-random cyclic binary signal specific to the satellite. Each satellite has its own pseudo-random code and its own carrier frequency. The pseudo-random code spreads out the satellite signal over a wide frequency band and embeds the satellite signal in the background noise. This limits the interferences between the satellite signals and the sensitivity to external perturbations.\nIn order to detect a satellite signal, the receiver generates a local code by means of a code generator driven by a digitally controlled oscillator, the local code reproducing the pseudo-random code of the satellite signal, and correlates the received signal with the local code.\nIn an initial acquisition phase, in which the receiver tries to detect the satellite signal, it is necessary to synchronize the locally generated local code with the received pseudo-random code of the satellite. To do this, the receiver comprises a code loop (or DLL, “Delay Lock Loop”) for subordinating the code oscillator. The code loop comprises a suitable code discriminator for receiving the result of the correlation, in order to determine a code error depending on the result of the correlation and for sending a corresponding code correction signal to the code oscillator.\nThe satellite navigation receiver receives the satellite signal in a direct line-of-sight and optionally, depending on the environment, in an indirect line-of-sight, for example after reflection on the ground, the sea or on buildings. These multiple paths perturb the detection of the satellite signal and are sources of error in calculating the positioning of the receiver.\nTo overcome this drawback, in the acquisition phase it is known to correlate the satellite signal with a combination of punctual, advanced and/or delayed local codes, so as to obtain a correlation function with which the pseudo-random code received in direct line-of-sight may be better discriminated.\nFR 2 739 695 discloses a Double Delta correlator which will be described in more detail hereafter. The theoretical correlation function of the Double Delta correlator has a narrow capture area surrounded by two so-called false lock-onareas, in which the correlation coefficient is zero.\nThe Double Delta correlator is efficient for an initial code error located in the capture area. Nevertheless, if the initial code error is located in a false lock-on area, the code loop operates in an open loop, which leads to a stationary measurement error.\nFurther, in a practical embodiment, the band pass of the receiver is not infinite so that the actual correlation function has secondary zeros in the false lock-on areas, over which the code loop may be locked on.\nIn order to overcome this drawback, FR 2 974 914 proposes a receiver comprising a Double Delta correlator associated with a false lock-on detector.\nNevertheless, this requires adaptation of the receiver for implementing the false lock-on detector, which is not necessarily possible upon updating a receiver."}
{"text": "Despite the functional safety of this system, it can nevertheless not be excluded that its functioning may be endangered by damage, especially in the case of vandalism, due to it necessarily requiring reflectors. In addition to this the expense in interwiring the light barriers is excessive."}
{"text": "The product described in JP-A-2001-326316, for example, exists as a conventional lead frame for resin-sealed semiconductor devices. In conventional lead frames such as this one, a large number of terminals are arranged around each of die pads, with tie bars being arranged in grid form on a plane in order to interconnect the large number of terminals to hanging leads.\nMeanwhile, illumination devices that each use light-emitting diode (LED) elements as a light source are used in recent years in general illumination, vehicle-mounted illumination, and displays, as well as in state indicators of various electrical household appliances, office automation machines and apparatuses, and vehicular devices. Some of these types of illumination devices include a semiconductor device fabricated by mounting LED elements in a lead frame."}
{"text": "1. Field of the Invention\nThe present disclosure relates to an image forming apparatus that a sheet material.\n2. Description of the Related Art\nA small and inexpensive stepping motor is often used as a drive source for a sheet feeding/conveying system in an image forming apparatus. The stepping motor is generally driven by a constant current chopper control system. While the stepping motor can contribute to achieving a compact and inexpensive structure, the stepping motor frequently causes a step-out phenomenon in which the rotation of a rotor cannot be synchronized with the input of a pulse signal. The step-out phenomenon occurs, for example, in an overload state to the pulse rate of the pulse output to the stepping motor from a drive circuit.\nAn image forming apparatus that performs multifarious types of image formation needs to support various kinds of sheets, such as plain paper and thick paper. Therefore, the required torque of the stepping motor varies significantly depending on the kind of sheets in use in some cases. With regard to the torque for causing a sheet to enter between conveying rollers made of a sponge material and disposed along a sheet conveying path, for example, the torque for thick paper (200 g/cm) becomes twice as high as the torque for plain paper (80 g/cm) in some cases. The torque of the stepping motor is determined by the value of the drive current. Therefor, the selection of the stepping motor and the selection of the drive current that determines the torque are determined on the assumption of using thick paper which faces severer conditions.\nHowever, plain paper is predominantly used in general offices or the like. Therefore, the constant use of the drive current set for thick paper causes more than necessary power consumption and noise generation. Thus, a technique to solve this problem is proposed (see Japanese Patent Application Laid-Open No. 2001-322734). According to the technique described in Japanese Patent Application Laid-Open No. 2001-322734, the drive current value is set for plain paper in a normal operation mode, and is set for thick paper only when a user operates to set thick paper.\nRecently, there is a proposal of a technique of optimizing the drive current based on the relationship between the maximum output torque of a stepping motor and the level of the drive current (see US 2011/0229235). The technique disclosed in US 2011/0229235 calculates an estimated load torque at the time of conveying a first sheet. The estimated load torque is used to determine a current value corresponding to the level of the target load torque at the time of conveying second and subsequent sheets.\nThe technique disclosed in Japanese Patent Application Laid-Open No. 2001-322734 is effective in preventing a step-out phenomenon. However, this technique still faces the problem when actual sheets in use differ from the set sheets, for example, when a sheet cassette for plain paper is removed and thick paper is fed in this sheet cassette.\nThe technique disclosed in US 2011/0229235 independently controls a plurality of stepping motors, thus complicating the control circuit when a plurality of stepping motors is used."}
{"text": "A user interface (UI), in the industrial design field of human-machine interaction, is the space where interactions between humans and machines occur. The goal of the interactions between a human, often referred to as a “user,” and a machine at the user interface is the user's control of the machine and its operations (e.g., through user input) and machine feedback (e.g., through program output). A graphical user interface (GUI) is a type of user interface that allows users to interact with software applications executing on electronic or computing devices through multimedia objects (e.g., images, videos, audios, etc.) rather than purely text commands. A user interface may allow users to interact with a social-networking system.\nA social-networking system, which may include a social-networking website, may enable its users (such as persons or organizations) to interact with it and with each other through it. The social-networking system may, with input from a user, create and store in the social-networking system a user profile associated with the user. The user profile may include demographic information, communication-channel information, and information on personal interests of the user. The social-networking system may also, with input from a user, create and store a record of relationships of the user with other users of the social-networking system, as well as provide services (e.g., wall posts, photo-sharing, event organization, messaging, games, or advertisements) to facilitate social interaction between or among users.\nThe social-networking system may send over one or more networks content or messages related to its services to a mobile or other computing device of a user. A user may also install software applications on a mobile or other computing device of the user for accessing a user profile of the user and other data within the social-networking system. The social-networking system may generate a personalized set of content objects to display to a user, such as a newsfeed of aggregated stories of other users connected to the user.\nA mobile computing device—such as a smartphone, tablet computer, or laptop computer—may include functionality for determining its location, direction, or orientation, such as a GPS receiver, compass, gyroscope, or accelerometer. Such a device may also include functionality for wireless communication, such as BLUETOOTH communication, near-field communication (NFC), or infrared (IR) communication or communication with a wireless local area networks (WLANs) or cellular-telephone network. Such a device may also include one or more cameras, scanners, touchscreens, microphones, or speakers. Mobile computing devices may also execute software applications, such as games, web browsers, or social-networking applications. With social-networking applications, users may connect, communicate, and share information with other users in their social networks."}
{"text": "The present invention relates to computing environments, and deals more particularly with techniques for obscuring information in messages to be exchanged over a communications network. In one aspect, the information comprises path name information and parameters for use in a Uniform Resource Locator (“URL”). In another aspect, the information comprises parameters used in forms.\nMessages exchanged over a communications network such as the Internet commonly contain a URL. As is well known, a URL is a subset of a Uniform Resource Identifier (“URI”) that specifies where a particular resource is located and a mechanism for retrieving it. The format of a URL is well known, and typically includes a host name and domain portion followed by a path name portion. A so-called “dynamic” form of URL is known, where the URL format provides for passing a string of one or more parameters, where each parameter comprises a name/value pair. The presence of the parameter string is denoted in the URL by specifying a question mark character, and the parameter string then follows this character. Within the parameter string, the name/value pairs are separated from one another by an ampersand character and an equal sign is used to separate the parameter name from the parameter value in each name/value pair.\nThere have been multiple cases of hacking web sites by modifying dynamic URL parameters, allowing the hackers to discover private data, corrupt the back-end application, and/or corrupt the data used by the back-end application. In one known security breach, for example, a URL included a parameter name/value pair for specifying a user's key for registered users of a web site to thereby pass this key value to the back-end application. Hackers discovered that by modifying the user key value, they were able to view name and address information of the user associated with the altered key value."}
{"text": "Statement of the Technical Field\nThe inventive arrangements relate to haptic interfaces, and more particularly to compact haptic interfaces which are designed to integrate with a primary controller.\nDescription of the Related Art\nRemote controlled unmanned vehicles are increasingly being used in a wide variety of robot applications such as explosive ordinance disposal, search and rescue operations, undersea salvage, and oil rig inspection/maintenance. As interest grows in robotic systems, providers are seeking to add haptic (force feedback) capability to their controllers. In many systems, a basic laptop-style controller already exists but these systems do not offer haptic feedback. Accordingly, there is a need for a haptic controller that can be used in connection with existing laptop-style controllers.\nIn many scenarios in which robots are used, conventional haptic interfaces are not well suited. These conventional haptic interfaces often have a form factor which lacks compactness and therefore do not work well. For example, conventional haptic interfaces are often designed for desktop consumer usage as opposed to mobile or portable robot operations. As such, these existing systems tend to be too large or have a form factors that makes them impractical for many applications."}
{"text": "Typically, in multi-processing computing systems, non-volatile memory/storage, “NVM”, such as fast Flash and Storage Class Memory is packed/interfaced as fast hard disk drive (SSD).\nFurther, some multi-processing and parallel computing systems currently implement networking standards or protocols, to maintain connections between multiple computer nodes (e.g., high performance computing nodes) and I/O nodes such as storage devices, e.g., hard disk drive memory storage devices. The Internet protocol suite based on the TCP and IP protocols is just one example. Other examples for networked storage device I/O connection standards for accessing hard disk storage devices and SSDs include the iSCSI protocol, which is an upper layer protocol based on TCP/IP. Running these protocols to exchange data between computing systems or computing systems and storage systems typically involves overhead due to copying the data to be communicated between involved application programs and the network protocol stack, and within the protocol stack.\nRDMA (Remote Direct Memory Access) is a communication paradigm to overcome this performance problem while transferring content of local memory to a peer hosts remote memory or vice-versa without involving either one's operating system during the actual data transfer and thus avoiding any data copy operation otherwise needed. Several protocol suites exist to implement an RDMA communication stack. Infiniband® (Trademark of System I/O, Inc., Beaverton, Oreg.), iWarp, and RoCEE (RDMA over Converged Enhanced Ethernet) are three example network technologies which can be deployed to implement an RDMA stack. These technologies use different network link technologies and different network packet formats to exchange RDMA messages between hosts.\nFurther, there currently exist switched fabric technologies infrastructure for server and storage connectivity such as the OpenFabrics Enterprise Distribution (OFED™ (Trademark of Open Fabrics Alliance, INC. California). OpenFabrics is an industry standard framework for a host implementation of the RDMA communication paradigm comprising the definition of an application programming interface (RDMA ‘verbs’ API) and generic user level and operating system level components to which network technology specific and vendor specific components can be attached in a standardized way. OpenFabrics is open-source software for RDMA and kernel bypass applications for use in high performance, highly efficient networks, storage connectivity and parallel computing. The OFED™ programming interface allows an application to access memory of a remote machine via RDMA directives such as RDMA Read, RDMA Write, RDMA Send and RDMA Receive.\nFIG. 1 particularly shows an example prior art model 10 implementing OFED™ connectivity standard used in context of RDMA in a network 12 having nodes that execute tasks and may access other networked nodes. Particularly, in the embodiment depicted in FIG. 1, host devices, e.g., networked computing devices, i.e., “peerA” device 16, and “PeerB” device 18, etc. on two different machines, and the communications network 12 that interconnects them create RDMA connections by first registering their respective virtual memory regions 26, 28. Then the peers 16,18, according to OFED™ standard directives, may perform RDMA READ/WRITE operations directly into a remote peer's virtual memory address space. For example, peer 16 may access virtual memory region 28 of peer 18 device via RDMA commands 27, and peer 18 may access virtual memory region 26 of peer 16 device via RDMA commands 25. The respective Infiniband adapters 22, 24 each performs the RDMA operations responsive to commands from the remote peer CPU and control path elements. The work queues 17 in PeerA and work queues 19 in PeerB have entries such as RDMA read or write directives and a respective control paths 13, 15 provides connectivity to trigger the respective Infiniband adapter to perform the RDMA.\nFurther, there currently exists a NVMe (Non-volatile memory Express) (www.nvmexpress.org/) describing a new standard to access PCI-attached NVM SSD's. This standard is based on an asynchronous multi-queue model, however, is still block based accessed (e.g., in a multiple byte unit such as 512 bytes, 4096 bytes, 16 Kilobytes etc.). That is, access to the fast Flash and Storage Class Memory (NVM) as persistent memory/storage is slowed down by classic “block access” methods developed for mechanical media (e.g., hard disks) in currently existing systems. This is a problem in that implementing block access methods increases NVM memory access and storage times.\nWith more particularity, current host controller devices, such as device 35 of FIG. 2, provide interfaces such as AHCI (HBA), SCSIe, or NVMe for PCI Express attached SSDs. In the case of NVMe, the controller 35 includes functionality for partitioning w/multiple ports, providing parallel IO, Scatter-gather-support, and up to 64,000 I/O Command (Send) Queues, with a maximum command depth of 64,000 entries, and Completion Queues.\nFIG. 2 shows an example of the NVMe standard block device access for a multi-core (multiprocessing) system 30 which interprets memory as a hard disk. It does not provide byte addressability. In FIG. 2, a user application via a controller interface 35 creates a work request(s) e.g., write, write a block of data to a specified location, e.g., a sector x, and fetch data for user at an address “y”. This work request(s) is (are) placed in a respective command queue or send queue (SQ) such as SQ 31 associated with a first computing device (e.g., processor Core, and SQ 34 associated with another computing device (e.g., processor Core n). As shown, each core 31, 34 of the multi-core (multiprocessing) system is in communication with the controller device 35 that is configured to process the queue entries for each host or core. The notification(s) are received at a respective completion queues (CQ) 33, 36 of each respective core, and the application reads/processes the CQ to know when the data request is completed.\nAs shown in FIG. 3, an OFED™ software stack framework 50 includes kernel-level drivers 52, channel-oriented RDMA and send/receive operations, kernel bypasses of the operating system 53, both kernel level application programming interface (API) 55 and user-level API 57 and providing services for parallel message passing (MPI), sockets data exchange (e.g., RDS, SDP), NAS and SAN storage (e.g. iSER, NFS-RDMA, SRP) and file system/database systems. The network and fabric technologies that provide RDMA performance with OFED™ include: legacy 10 Gigabit Ethernet, iWARP, RoCE and 10/20/40 Gigabit InfiniBand.\nThe OFED™ framework defines access to remote memory at byte granularity and thus avoids the drawbacks of block-based access of protocols such as NVMe. Nevertheless, the OFED™ framework is currently defined for only accessing remote computer memory via a network link, and thus cannot be used to access local Non Volatile Memory.\nA new byte oriented access method for local NVM is necessary. This access method must support highly parallel or multicore systems."}
{"text": "Since the development of computer networks, security has been a concern of administrators of computers and computer networks. As a result, many different methods of securing computers have been proposed.\nOne such security method is a firewall. Firewalls provide for security of computers by regulating what data is allowed into and out of a computer or a computer network to which the computer is connected.\nAnother method for computer network security is Internet Protocol Security (IPsec), also called connection security. IPsec is a protocol that secures connections between two computers, or a computer and a computer network having a device supporting IPsec."}
{"text": "The recycling of old rubber tires presents a very difficult problem. In the first place, there is considerable difficulty in separating embedded metal parts from the rubber of the tire before the tire runs through the pulverizing machine. Moreover, there is a possibility of reusing the rubber only when the rubber is reduced practically to powder form because it is only in this form that old rubber can be used as an additive in the rubber industry for the production of new rubber ware.\nIn order to comminute old rubber tires, cutting and shredding tools have been used. According to DE 29 11 251 C2, the tire is clamped in a clamping device, with three outwardly slidable pins, whereupon a tangential cut is made with rotating disk knives to remove the tread with the side strips from the metal inlay containing beads and subsequently other circular disk-form knives are used to cut the rubber into strips which are then fed into cross cutting apparatus. The rubber pieces thus obtained required further comminuting.\nOther old rubber tire disintegrating machines are shredders which, with rotating knife cylinders, attack the circumference of the tread and comminute the tread as well as the side strips to chips (SU 36 93 894 A1). A machine according to SU 13 88 294 A1 works in a similar manner. Here the tire is gripped on both sides by saucers and is stressed and pressed so that the tread is bent and deformed and in this narrowed form is subjected to the knives of a shredding cylinder. Comminuting is effected in a similar manner by the comminuting machine of DE 37 04 725 A1, in which the comminuting tools arranged on the rotating comminuting cylinder are small, hard plates, and the comminuting tools themselves are arranged in a particular form.\nThe old rubber pieces obtained from all of these comminuting processes are technically not yet usable in the rubber industry as an additive. The rubber pieces thus obtained are so large, of such a dissimilar shape and so nonhomogeneous that further working and comminuting presents great difficulty.\nRubber scrap is also produced in the production of new rubber ware when, through production error or mixing error or through use of unsuitable mixtures, damaged goods are produced.\nApparatus for further working such rubber scrap is known through DD 265 855 A1. With this known apparatus, an extruder with several work zones is connected by a radial discharge pipe with a mill having grinding disks defining a conical grinding gap. Different work zones are formed in the extruder through different formation of different screw sections. Bringing the compressed rubber scrap from the extruder into the grinding space is difficult and the grinding of the rubber is problematic.\nFor rubber is an extremely difficult to comminute material, which is in many respects nonhomogeneous. Through different exposure to the sun's rays, and different loading during the life of an automobile tire, as well as different life spans, the old rubber scrap to be worked on is differently aged. The pieces of old rubber scrap have strongly aged portions, which are relatively easy to pulverize. Other, less aged parts, on the contrary, still have a high degree of elasticity and can hence be comminuted by grinding only with great difficulty or not at all. In most cases, strongly aged portions are integral with less aged portions. They thereby present particular difficulty in comminuting. This difficulty is increased by the presence of foreign material in the form of metal parts in the rubber scrap which are not removable prior to grinding because they are in part surrounded by the rubber and in part the rubber is vulcanized to them. These metal parts have a destructive effect on the grinding disks."}
{"text": "The present invention relates to a mobile communications device and methods of operation. More specifically, embodiments of the present invention relate to a mobile communications device, such as a smart phone, including an ultraviolet light source, and methods of controlling the ultraviolet light source using the smart phone.\nThe inventor of the present invention is aware of the use of ultraviolet light for disinfectant purposes. Currently, there are few stand-alone products on the market that provide ultraviolet light for cleaning surfaces or purifying water. One such product is a hand held UV wand that is plugged into a wall socket, and waved over surfaces; and another such product is a hand-held unit that runs on batteries, and is inserted to sanitize a bottle of water.\nSome drawbacks contemplated by the inventor, to such devices include the high power consumption of such devices limit utility of such devices. For example, surface sanitizers are typically bulky and need to be powered by plugging them into a wall socket; and portable water sanitizers use batteries, but drain them quickly. Additional drawbacks are when the user travels, the user must remember to bring along. Because of gadget overload, such dedicated ultraviolet light (UV) sources are not believed to be widely adopted.\nIt is desired to have an ultraviolet light source without the drawbacks described above."}
{"text": "The fabrication of integrated circuits (ICs) involves the formation of features on a substrate that make up circuit components, such as transistors, resistors and capacitors. The devices are interconnected, enabling the IC to perform the desired functions. Interconnections are formed by forming contacts and conductive lines in a dielectric layer using, for example, damascene techniques. A damascene structure, for example, includes a via or contact hole in a lower portion and a trench in an upper portion. The via serves as a contact to a device while the trench contains the conductive line for connecting the device to, for example, other devices.\nIt is important for via contacts in an upper interconnect level to be coupled properly with the underlying conductive lines of a lower interconnect level. However, lithographic errors, such as via misalignment or develop critical dimension (DCD) which is larger than the desired CD, may result in via contacts being extended beyond or displaced from the underlying conductive lines. As a result, the separation distance between the via contact and the adjacent lower conductive line is reduced. This may lead to electrical shorts between the via contact above and the adjacent conductive line below. Moreover, via contacts which extend beyond or fails to land on the underlying conductive lines may cause electrical shorts to the layer beneath the lower conductive layer due to punch through. These phenomena adversely render the IC malfunction.\nFrom the foregoing discussion, it is desirable to provide a device which is devoid of the above-mentioned problem, thus increasing the reliability of the IC. It is also desirable to provide a cost effective process for forming the device."}
{"text": "With the development of network technologies, there is an increasing quantity of users participating in various social activities by using networks. To enrich lives, some users record, by using a recording device, a video in which they play a role, upload the recorded video to a network, wait for another user to record a video of a partner role, and synthesize a costarring video from the video of the self-played role and the video of the partner role.\nHowever, in the conventional manner, a role in a video played by a user himself is recorded first, and a remaining role in the video is filled with a piece of still picture or an animation, waiting to be played by someone actively. The efficiency of costarring is greatly reduced when the waiting time is excessively long."}
{"text": "Prediction plays a critical role in video coding, where current video data is coded by using previously reconstructed video data or other derived data as predictor. When coding or decoding a block, a prediction block corresponding to intra-prediction or inter-prediction is generated first. The prediction residue is then formed at the encoder side based on the current block and the prediction block. At the decoder side, the prediction block is used to reconstruct the reconstructed samples using the coded prediction residue.\nBesides inter-prediction and intra-prediction, a new prediction mode named intra-block copy (IBC) has been adopted in the screen content coding (SCC) profile for the high efficiency video coding (HEVC) standard. IBC can take advantage of copying contents in the same picture. FIG. 1 illustrates an example of IBC, where a reference block (120) in the current picture is copied to the current block (110) as the prediction if IBC is applied. The reference block is located by applying a block-copying vector (BV, 130). The example in FIG. 1 shows that the BV points to the reference block horizontally. However, the BV can be a two-dimensional vector not restricted to be horizontal. The samples in the reference block have been reconstructed already before the current block is coded or decoded.\nInter simplified depth coding (InterSDC) is adopted into 3D-HEVC (three-dimensional HEVC) as a special prediction mode for depth coding. When InterSDC is used, a normal inter-prediction is first performed for the current block. A coded offset is then added to each pixels in the prediction block. If Pi,j represents the prediction value at pixel position (i, j) after performing the normal inter-prediction and Offset is the coded offset for this block, the final prediction value at pixel position (i, j) is Pi,j+Offset. With the InterSDC mode, no residues are coded. Thus the final prediction value will be output as the reconstructed value.\nDepth lookup table (DLT) is adopted into 3D-HEVC. Since there are often only a small number of pixel values appearing in the depth component, the DLT signals those valid values from the encoder to the decoder. When a coding unit (CU) is coded in the intra simplified depth coding (SDC) mode or depth map modeling (DMM) mode, DLT is used to map the valid depth values to DLT indices. The DLT index is more efficient for compression. FIG. 2 illustrates an example of the DLT process, where only five pixel values (i.e., 50, 108, 110, 112 and 200) among the 256 possible pixel values appear in a picture block (e.g., a coding unit). The five pixel values can be represented using five DLT indices (i.e., 0, 1, 2, 3 and 4). The DLT can be signaled in picture parameter set (PPS). And it is up to the encoder to derive the DLT during the encoding process.\nSince prediction values come from reconstructed pixels, there are distortions between the prediction values and the original values even if the original pixels in the current block and the original pixels in the reference block are exactly the same. Furthermore, since the reconstructed signal loses high-frequency information generally, the quality of prediction is deteriorated more when there are sharp pixel value changes in the reference block. FIG. 3 and FIG. 4 demonstrate two examples, where there are two and three segments with sharp edges in a block.\nAccordingly, it is desirable to develop a method for prediction that can take into consideration of sharp transitions in the prediction block to improve coding performance."}
{"text": "The invention relates to a process for the preparation of mono- and oligoisocyanates by reacting primary amines with phosgene in the presence of a catalyst.\nIsocyanates are industrial products which have a large number of uses in the field of polyurethane plastics. However, certain isocyanates are also used in the preparation of pharmaceutical active ingredients.\nThe synthesis of isocyanates by reacting amines with phosgene has been known for some time. In principal, two processes are described in the literature, one of which is carried out at atmospheric pressure and the other is carried out at increased pressure. Phosgenation under increased pressure is disadvantageous since it requires much more complex industrial apparatus to control the increased safety risk, the release of phosgene.\nFor sulfonyl isocyanates, U.S. Pat. No. 3,371,114 and U.S. Pat. No. 3,484,466 disclose a preparation process at atmospheric pressure in which a solution of a sulfonylamide and an isocyanate as catalyst in an inert solvent is reacted with phosgene. In the process, the corresponding sulfonylurea is formed as an intermediate, which reacts with phosgene to give the desired sulfonyl isocyanate.\nAlkyl and aryl isocyanates are usually prepared by the phosgenation process, described, for example, in Houben-Weyl, Methoden der organischen Chemie [Methods in Organic Chemistry], 4th Edition, Volume E4, pages 741-751, Georg Thieme Verlag Stuttgart, 1983, from the corresponding amines in two phases at atmospheric pressure. In the first phase, the cold phosgenation, the amine is reacted with an excess of phosgene in very dilute solution and at low temperatures to give the corresponding carbamyl chloride, from which, in the second phase at elevated temperature, the hot phosgenation, the isocyanate forms. Aliphatic and cycloaliphatic primary amines are more difficult to phosgenate because of their increased basicity compared with aromatic amines, and lead to an increased formation of byproducts. A disadvantage of these processes is, in addition to the fact that the phosgenation is carried out in two phases, the formation of an intermediate solids suspension of sparingly soluble carbamyl chloride and amine hydrochloride, which in turn renders an increased dilution of the reaction medium necessary in order to prevent deposits and blockages of parts of the equipment. Because of the accumulation of solids which occurs, this process cannot be carried out continuously at atmospheric pressure. Furthermore, symmetrically N,Nxe2x80x2-substituted urea forms as a byproduct, the formation of which can only be suppressed at the expense of drastically reduced space-time yields.\nAliphatic and cycloaliphatic amines are frequently used in the form of their salts in the cold/hot phosgenation. However, these salts are sparingly soluble in the reaction medium, meaning that additional reaction stages and very long reaction times are necessary.\nFurthermore, it is known from GB 1 114 085, U.S. Pat. No. 3,492,331 and H. Ulrich, Chemistry and Technology of Isocyanates, Wiley and Sons, 1996, pages 328-330, to optimize the reaction of primary amines with phosgene by the addition of catalysts such as dimethylformamide, phenyltetramethylguanidine, 2,4,6-trimethylpyridine or carbodiimidazole. Some of these catalysts must be used in equimolar amounts and form sparingly soluble salts under the reaction conditions.\nIt is an object of the present invention to provide a process, which can be used for aliphatic, cycloaliphatic, araliphatic and aromatic primary amines, for the preparation of the corresponding mono- and oligoisocyanates using phosgene which can be carried out either continuously or batchwise at atmospheric pressure, does not have the above disadvantages and provides the corresponding mono- and oligoisocyanates in good yields and high selectivities.\nWe have found that this object is achieved by a process for the preparation of aliphatic, cycloaliphatic, araliphatic and aromatic mono- and oligoisocyanates by phosgenation of the corresponding primary amines at atmospheric pressure, in which\na) a catalytic amount of a monoisocyanate (isocyanate a) is introduced into an inert solvent together with phosgene,\nb) the primary amine is added, and\nc) the resulting reaction mixture is reacted with phosgene.\nThe net equation underlying the process is given in scheme 1 below. \nThe process according to the invention has the advantage that it can be used for a large number of amines. The phosgenation is carried out according to the invention, while avoiding a division into cold and hot phosgenation, in a narrow temperature interval and at atmospheric pressure, the intermediate formation of sparingly soluble suspensions being avoided. The desired isocyanate is formed in the process, with complete conversion of the amine, in high yields and high selectivity in significantly shortened reaction times without symmetrically substituted N,Nxe2x80x2-urea being formed from the amine as a byproduct. Since the formation of urea from the amine is not observed, it is possible by means of the process according to the invention to significantly increase the concentration of the amine in the reaction solution and thus the space-time yields. In addition, it is advantageous that the process according to the invention can be carried out either batchwise or continuously since there is no accumulation of solids.\nThe process according to the invention gives aliphatic, cycloaliphatic, araliphatic and aromatic mono- and oligoisocyanates of the formula I\nR1xe2x88x92N=C=Oxe2x80x83xe2x80x83(I)\nThe radical R1 in formula I corresponds to the radical R1 in formula IV of the amines used in the process according to the invention, which are discussed later and to which reference is made here. Preference is given to preparing mono- and diisocyanates by the process of the invention. Of lesser importance in practice, but preparable in principle are isocyanates having 3 and more isocyanate groups.\nThe sole catalyst used is a monoisocyanate of the formula II (isocyanate a)\nR2xe2x88x92N=C=Oxe2x80x83xe2x80x83(II)\nor mixtures thereof, in which R2 is aliphatic, cycloaliphatic, aromatic or araliphatic radicals. These can be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. The radical R2 must, however, be inert toward phosgene, thus excluding radicals which carry NH, OH and SH groups. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 3 to 30 carbon atoms, preferably 3 to 10 carbon atoms. Examples of aliphatic radicals are methyl, ethyl, propyl, n-butyl, isobutyl and sec-butyl.\nSuitable cycloaliphatic radicals are those which have 3 to 20 carbon atoms, preferably 3 to 10 carbon atoms, such as, for example, cyclopentyl and cyclohexyl.\nThe aromatic radicals can be unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. Preference is given to the aromatic radicals which are mono- or disubstituted. Examples of aromatic radicals are phenyl, chlorophenyl, o-, m- and p-tolyl.\nSuitable araliphatic radicals are radicals having 7 to 12 carbon atoms, such as, for example, benzyl, although preference is given to a radical of the formula III\n\nin which R3 and R4 can be identical or different and can be hydrogen or an aliphatic, cycloaliphatic, araliphatic or aromatic radical, in which case these may be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. The radicals R3 and R4 must, however, be inert toward phosgene, thus excluding radicals carrying NH, OH and SH groups. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 1 to 30 carbon atoms, preferably 1 to 10 carbon atoms. Examples of aliphatic radicals are methyl, ethyl, propyl, n-butyl, isobutyl and sec-butyl.\nSuitable cycloaliphatic radicals are those which have 3 to 20 carbon atoms, preferably 3 to 10 carbon atoms, such as, for example, cyclopentyl and cyclohexyl.\nThe aromatic radicals can be unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. Preference is given to the aromatic radicals which are mono- or disubstituted. Examples of aromatic radicals are phenyl, chlorophenyl, o-, m- and p-tolyl. R5 is fluorine, chlorine, bromine or a C1- to C4-alkyl chain which may optionally be interrupted by a heteroatom, for example oxygen or sulfur. Preference is given to the monoisocyanate (isocyanate a), a linear aliphatic alkyl isocyanate, particularly preferably a linear aliphatic alkyl isocyanate having 3 to 10 carbon atoms, where the alkyl chain may optionally be branched, particularly preferably n-butyl isocyanate. In the case of the preparation of monoisocyanates by the process according to the invention, it is also possible to use the monoisocyanate desired as product in catalytic amounts as isocyanate a. It is, however, preferable to use an isocyanate a which is not identical to the product (isocyanate of the formula (I)) as catalyst.\nThe isocyanate a is generally used in a catalytic amount of from 0.01 to 25 mol %, preferably 0.5 to 20 mol %, particularly preferably 1 to 15 mol %, based on the amine.\nIn step a) of the process according to the invention, the isocyanate a is introduced into a solvent which is inert toward phosgene. Preferred solvents are hydrocarbons. Particular preference is given to mono- or polysubstituted aromatic hydrocarbons, such as toluene, o-, m- or p-xylene, ethylbenzene, chlorobenzene or 1,2-, 1,3- or 1,4-dichlorobenzene or mixtures thereof. Particular preference is given to xylenes, chlorobenzene and dichlorobenzenes as solvent.\nThe process according to the invention is carried out at atmospheric pressure and generally at a reaction temperature, which is largely constant over the three steps of the process a), b) and c), of from 20 to 200xc2x0 C., preferably from 20 to 150xc2x0 C., particularly preferably from 50 to 120xc2x0 C.\nThe monoisocyanate (isocyanate a) dissolved in the inert solvent is, in step a) of the process according to the invention, heated to the reaction temperature and mixed with an excess of phosgene, based on the monoisocyanate. The molar ratio between phosgene and isocyanate a is preferably 100:1 to 5:1. Particular preference is given to using phosgene in an excess of 20:1 to 5:1, based on isocyanate a.\nThe phosgene can be used in the process according to the invention in steps a) and c) in gaseous form or in condensed form, preference being given to introducing the phosgene in gaseous form.\nIn the process according to the invention, suitable primary amines are amines of the formula IV\nR1xe2x88x92NH2xe2x80x83xe2x80x83(IV)\nwhich are obtainable by methods known from the literature or as a commercially available product.\nR1 stands for aliphatic, cycloaliphatic, aromatic or araliphatic radicals. These can be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. Radicals which contain OH and SH groups are excluded. Substitution by amino groups, on the other hand, is possible and leads to di- or oligoisocyanates, preference being given to diisocyanates. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 3 to 30 carbon atoms, preferably 3 to 10 carbon atoms. Examples of aliphatic radicals are propyl, butyl, pentyl and hexyl. Suitable cycloaliphatic radicals are those having 3 to 20 carbon atoms, preferably those having 3 to 10 carbon atoms, such as cyclopentyl and cyclohexyl.\nThe aromatic radicals can have 4 to 10 carbon atoms, can be heteroaromatic radicals, and unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. The aromatic radicals are preferably mono- or disubstituted. Examples of aromatic radicals are phenyl and chlorophenyl, o-, m-, p-tolyl, furfuryl, thiophenyl and pyrrolyl.\nSuitable araliphatic radicals are radicals having 7 to 12 carbon atoms, such as, for example, benzyl, although preference is given to a radical of the formula V \nin which R6 and R7 can be identical or different and can be hydrogen or an aliphatic, cycloaliphatic or aromatic radical, in which case these can be substituted by heteroatoms, or their carbon chains can be interrupted by heteroatoms, such as oxygen and sulfur. The radicals R6 and R7 must, however, be inert toward phosgene, meaning that radicals carrying NH, OH and SH groups are excluded. The aliphatic radicals can be arbitrarily branched or unbranched, saturated or unsaturated. They contain 1 to 30 carbon atoms, preferably 1 to 10 carbon atoms. Examples of aliphatic radicals are methyl, ethyl, propyl, n-butyl, isobutyl and sec-butyl.\nSuitable cycloaliphatic radicals are those which have 3 to 20 carbon atoms, preferably 3 to 10 carbon atoms such as, for example, cyclopentyl and cyclohexyl.\nThe aromatic radicals can be unsubstituted or arbitrarily substituted by alkyl or aryl substituents or heteroatoms. Preference is given to the aromatic radicals which are mono- or disubstituted. Examples of aromatic radicals are phenyl, chlorophenyl, o-, m- and p-tolyl. R8 is a halogen atom (F, Cl, Br) or a C1- to C4-alkyl chain, which may optionally be interrupted by a heteroatom, for example oxygen or sulfur. Primary amines of the formula IV which can be used are also amines in enantiomerically pure, optically active form or mixtures of the enantiomers. The process according to the invention advantageously leads during the reaction and the work-up to only a low degree of racemization, which is generally below 2%.\nIn the process according to the invention, preference is given to using cycloaliphatic, araliphatic and aromatic amines, in particular those which are further processed in industrial processes, such as, for example, hexamethylenediamines, cyclohexylamine, isophoronediamine, toluylenediamine, aniline, and optically active amines, such as, for example, R-(+)- and S-(xe2x88x92)-phenylethylamine.\nIn step b) of the process according to the invention, the primary amine is added to the solution of the monoisocyanate (isocyanate a), saturated with phosgene, preferably in an inert solvent. This can be carried out, for example, in the presence of a stream of inert gas, for example nitrogen. Preferred solvents are the solvents used in step a) of the process according to the invention. Particular preference is given to using the same solvent as in step a). The resulting reaction mixture is then reacted with phosgene in step c) of the process according to the invention.\nThe molar ratio between primary amine and the total amount of the amount of phosgene introduced in steps a) and c) is 1:10 to 1:1, particular preference being given to introducing phosgene in an amount of from 120 to 180 mol %, based on the amine.\nThe process according to the invention can in principle be carried out both continuously and batchwise, preference being given to a continuous operation. The process can be carried out in any desired apparatus which is suitable for a reaction with phosgene. For example, a phosgenation apparatus is used which comprises an attached carbonic acid condenser or a downstream battery of high-efficiency condensers for the condensation of phosgene.\nFor the preferred continuous procedure, a phosgenation apparatus is used which comprises a main reactor and after-reactor, a stripping column and, connected downstream therefrom, a condenser. In the case of the continuous procedure, a mixture of the desired isocyanate of the formula I, the catalytic amounts of the isocyanate of the formula II (isocyanate a) and phosgene in the inert solvent is preferably initially introduced.\nAfter an after-reaction, excess phosgene and solvent are generally stripped out at temperatures from 30 to 80xc2x0 C., preferably 40xc2x0 C. to 60xc2x0 C., condensed in the condenser and returned together to the after-reactor.\nThe reaction product is worked up using methods known to the person skilled in the art. Preferably, the desired isocyanate of the formula I is isolated by fractional distillation. The isocyanate a used as catalyst is preferably recovered by distillation and can be reused for further batches. The desired mono- and oligoisocyanates can be obtained in the process in yields of generally 60 to 95%, based on the primary amine used.\nThe examples below illustrate the invention further."}
{"text": "1. Technical Field\nThe present invention relates to a robot system.\n2. Related Art\nIn the related art, a robot provided with a robot arm is known. In the case of the robot arm, a plurality of arms (arm members) are linked to each other via a joint portion, and, for example, a hand is mounted on the arm furthest on the tip end side (the most downstream side) as an end effector. The joint portion is driven by a motor, and the arm rotates by the driving of the joint portion. In addition, the robot grabs an object by the hand, moves the object to a predetermined position, and performs predetermined work, such as assembly.\nAs such a robot, in JP-A-2014-46401, a vertical articulated robot is disclosed. The robot described in JP-A-2014-46401 is configured to perform an operation of moving a hand with respect to a base to a position different by 180° around a first rotating axis that is a rotating axis (rotating axis that extends in a vertical direction) furthest on the base end side (most upstream side), by rotating a first arm which is the arm furthest on the base end side with respect to the base around the first rotating axis.\nIn the robot described in JP-A-2014-46401, when moving the hand with respect to the base to a position different by 180° around the first rotating axis, a large space for preventing interference of the robot is necessary."}
{"text": "Initiators are devices including a charge initiated to ignite or begin the burning of a larger main charge or propellant.\nStandards such as MIL-STD-1901A include a safety requirement to the effect that an initiator subject to 500 volts due to an errant voltage shall not detonate or deflagrate.\nAlthough various devices for protection against errant voltages are well known, to date those skilled in the art have failed to provide a suitable initiator which meets the 500 volt no deflagration safety requirement. Previous attempts at engineering such a suitable initiator resulted in undue complexity and initiators which are difficult to install or incorporate into existing systems."}
{"text": "Field\nThe present disclosure relates to a method for executing a menu in a mobile terminal for providing a convenient input manner and a mobile terminal using the same.\nBackground\nGenerally, terminals can be classified into mobile/portable terminals and stationary terminals. The mobile terminals can be further classified into handheld terminals and vehicle mount terminals according to possibility of user's direct portability.\nAs functions of the terminal are diversified, the terminal is implemented as a multimedia player provided with composite functions such as photographing of photos or moving pictures, playback of music or moving picture files, game play, broadcast reception and the like, for example.\nTo support and increase the terminal functions, the improvement of structural parts and/or software parts of the terminal may be considered.\nRecently, as the mobile terminal comes to be able to process the multi-task, the input manner which enables the user to execute a function while executing other function has been studied."}
{"text": "Polymer electrolyte fuel cells (or solid polymer fuel cells) have a basic component called a membrane electrode assembly (MEA). The MEA is obtainable by laminating an electrode membrane (a catalyst layer or an electrode catalyst membrane) on either side of a polymer electrolyte membrane and putting the resulting laminate between a fuel gas supply layer and an air supply layer; the electrode membrane contains, as a main component, a platinum group metal catalyst supported on a carbon powder, the polymer electrolyte membrane is an ion exchange membrane, and the fuel gas supply layer and the air supply layer are an electroconductive porous membrane. For the MEA, both electrolyte membrane and electrode membrane contain an ion exchange polymer. The electrolyte membrane and the electrode membrane are usually formed by a casting method and/or a coating method, or other methods. A method of laminating the electrolyte membrane and the electrode membrane usually comprises allowing an electrolyte membrane formed on a support to closely contact with an electrode membrane formed on a support, thermocompression bonding these membranes at about 130 to 150° C. (about 150 to 200° C. according to a material used) and a pressure of about 1 to 10 MPa, and then separating the supports from these membranes. Thus, as the support, a release film is used. The release film requires a moderate (or reasonable) releasability (peelability) from the electrolyte membrane and the electrode membrane and a moderate (or reasonable) adhesion to the electrolyte membrane and the electrode membrane. The ion exchange polymer contained in the electrolyte membrane and the electrode membrane has a unique structure that has a main chain being a highly releasable fluoropolymer and a side chain containing a poorly releasable sulfonic acid group. This makes it difficult to predict the releasability behavior and to balance the releasability and the adhesion. The release film is practically laminated on a substrate film having high mechanical characteristics in order to achieve improved handling or improved production. Unfortunately, it is difficult to improve the adhesion of the release film to a commonly-used substrate film having no reactive groups or other groups. The release film for the production of a fuel cell further requires heat resistance for reasons of production process. The release film, which is produced by a roll-to-roll processing in terms of efficient production, also requires flexibility. As the release film, a fluorine-containing film is widely used. The fluorine-containing film has excellent heat resistance, releasability, and stain resistance, while the film is expensive and less combustible in waste incineration after use, and easily generates a poisonous gas. The film, which has a low elastic modulus, is difficult to produce by a roll-to-roll processing. In place of the fluorine-containing film, a release film containing a cyclic olefin polymer is also reported.\nJapanese Patent Application Laid-Open Publication No. 2010-234570 (JP-2010-234570A, Patent Document 1) discloses a release film made of a cycloolefin copolymer and also describes a release film formed by coating a solution of a cycloolefin copolymer on a substrate film such as a poly(ethylene terephthalate) film. In working examples of the document, a solution containing a copolymer of ethylene and norbornene is cast on a poly(ethylene terephthalate) film using a flow-casting apparatus to form a release film having a thickness of 0.5 μm.\nThis release film has an excellent releasability from an electrolyte membrane or an electrode membrane containing an ion exchange polymer, while the release film has a low adhesion to a poly(ethylene terephthalate) film and is easily separated from the substrate film. This reduces the handleability of the release film or the production of the MEA. Further, the releasability of the release film from the electrolyte membrane or the electrode membrane is insufficient for a high-temperature production process of the MEA.\nJapanese Patent Application Laid-Open Publication No. 2009-102558 (JP-2009-102558A, Patent Document 2) discloses a laminated film stacked on a plastic substrate. The laminated film has a first layer composed of a chlorine-containing polymer, and a second layer laminated on the first layer and composed of a cyclic olefin polymer. Each of the first layer and the second layer is formed by coating. This document describes that the laminated film is utilizable for an industrial material such as an industrial release film, a packaging film for food, pharmaceutical, chemical, or other products, and an optical member such as a polarization plate for liquid crystal. The document also describes a copolymer of a cyclic olefin and a chain olefin as a preferred cyclic olefin polymer.\nHowever, this document is silent on a fuel cell. Further, the complicated structure of the laminated film would reduce fuel cell production or film handleability if the film is used as a release film for producing a fuel cell. This document fails to disclose the reason why a copolymer of a cyclic olefin and a chain olefin is preferred among cyclic olefin polymers.\nJapanese Patent Application Laid-Open Publication No. 2000-95957 (JP-2000-95957A, Patent Document 3) discloses a reversible thermotropic molding material having a first component and a second component; the first component consists of a transparent plastic material, and the second component consists of non-liquid crystal plastic material which is thermodynamically immiscible with the first component and is different in temperature dependency of refraction index from the first component. This document describes that an example of the first component includes a cycloolefin and that an example of a monomer constituting the second component includes vinylidene chloride. In working examples of this document, a molding material having a copolyamide in combination with a terpolymer containing a styrenic monomer or other monomers is prepared. The document also discloses that the applications of the molding material include reversible thermotropic plastic molding materials such as shielding systems for regulating the amount of light and temperature, for example, the glazing of buildings, greenhouses, cars, solar collecting systems and the like.\nThis document is also silent on a fuel cell. Further, this document fails to describe the meanings of combination of a cyclic olefin polymer and a chlorine-containing polymer."}
{"text": "In naval and other maritime vessels, such as aboard ships and submarines, there is a great need to deliver electrical power through spaced apart compartments, each of which can be independently sealed off in a water and air tight manner during emergency situations. Typically, an emergency situation aboard a naval vessel might be a fire in one or more compartments, or water leakage occurring therein, occasioned, for example, during battle.\nDuring such emergency conditions, it is imperative that each compartment in which the emergency condition arises, be sealed off and contained from the rest of the ship. To ensure this, a number of measures are conventionally undertaken.\nFirstly, naval power cable with water impenetrable blocking or filler material between the outer jacket and inner conductors, is generally used aboard such vessels. Secondly, electrical power cable extending through adjacent compartments is passed through water blocking cable transits, at which a water and air tight seal is sought between the transits and cable jacket. Thirdly, compartments aboard the vessel through which electrical power cable passes, are capable of being closed shut in a water and air tight manner during emergency situations.\nIn FIG. 1 of the drawings, a typical prior art multi-conductor power cable for use aboard naval vessels, is illustrated. In general, such prior art power cable 1 comprises a core conductor assembly 2 consisting of triplexed insulated conductors 3, a water blocking layer 4 extruded over the core conductor assembly, a non-uniform layer of binder tape 5 helically wound over the blocking layer, and a jacket layer 6 extruded over the binder tape, as shown. In order to install naval power cable through the various compartments within a vessel, the cable must permit bending at severe angles during cable routing operations.\nWhile it has been possible to satisfy the general requirements of naval power cable, prior art naval power cable of the type described above has nevertheless suffered from several significant shortcomings and drawbacks.\nIn order to expose the insulated electrical conductors for spreading and connection to electrical equipment in cable spreading rooms, large lengths of water blocking and jacket material must be stripped away to freely expose the core conductor assembly. With such prior art power cable, removal of the water blocking and jacket material requires cutting each layer away bit by bit using sharp cutting tools, the nature of which poses great safety hazards to shipyard workers. In addition to requiring arduous effort and numerous hours and sometimes a day or more to strip away the water blocking and jacket material, often the shipyard workers wielding the sharp cutting tools inadvertently cut into the core conductor assembly, damaging the electrical insulation on the conductors. Frequently, these workers also injure themselves during such stripping operations. For example, workers can cut themselves with the sharp cutting tools. Also they are exposed to the risk of slipping and falling on rubber chips and shavings which fall to the vessel floor during cable stripping operations.\nIn an effort to avoid damaging the electrical insulation during jacket and blocking material stripping operations, prior art power cable illustrated in FIG. 1 has required the overlapping layer of binder tape to help shipyard workers determine the location of the insulated electrical conductors and the blocking material. Consequently, prior art naval power cable has been generally characterized by surface irregularities in the outer jacket layer due to the underlying helically wound binder tape. These surface irregularities present difficulty in establishing high quality air and water seals desired at the interface of the cable outer jacket and cable transit sealing material. As a result, when a dozen or more of such prior art power cables pass through a cable transit, the risk associated with smoke and water leakage from one compartment to adjacent compartments, is greatly magnified during emergency situations.\nDue to the presence of the overlapping binder tape between the water blocking and outer jacket layer, wrinkles often form in the outer jacket of prior art cable as a result of cable bending motion naturally occasioned during installation (i.e., routing) of power cable over cable support hangers, and through conduits and cable transits. Consequently, with such wrinkles, the risk of damage to the outer jacket of prior art naval power cable is substantially increased during the cable installation process.\nThus, there is a great need in the art to provide improved naval power cable and a method of installing the same without the accompanying shortcomings and drawbacks of the prior art.\nAccordingly, it is a primary object of the present invention to provide a naval power cable which can be worked through severe bend angles during installation and yet be easily stripped of its water blocking material and outer jacket layer with a substantial reduction in stripping time and an added measure of safety.\nAnother object of the present invention is to provide such naval power cable, in which the inside surface of the outer jacket is more tenaciously bonded to the outside surface of the water blocking material than with the use of adhesives.\nAnother object of the present invention is to provide a naval power cable, in which the water blocking material is not internally damaged as the cable is being severely bent during cable working operations.\nAnother object of the present invention is to provide a naval power cable, in which the water blocking material and outer jacket layer are tenaciously bonded together and can be easily cut without the use of sharp cutting tools so that large lengths of power cable can be quickly and safely stripped aboard a naval vessel, by simply peeling away at most a few integrally bonded strips of water blocking and jacket material.\nAnother object of the present invention is to provide such naval power cable, in which, the outer jacket is essentially resistant to wrinkling when manipulated during the installation aboard a ship or like naval vessel.\nAnother object of the present invention is to provide such naval power cable, in which the outer jacket remains essentially uniform over its surface characteristics during installation and while embraced within the sealing material of cable transits, so as to provide improved water and air seals therebetween when subjected to hydrostatic pressures and smoke typically presented during emergency situations.\nYet a further object of the present invention is to provide a novel method of installing electrical power cable through isolatable compartments aboard a naval vessel, in a manner which reduces cable stripping time and provides an additional measure of shipyard worker safety.\nAn even further object of the present invention is to provide a simpler method of manufacturing such shipboard power cable.\nThese and other objects will become apparent hereinafter and in the claims."}
{"text": "In the related art, games are known in which an image of a part of a user's body such as a head is captured by a video camera, a predetermined area (e.g., an area including an image of an eye, a mouth, a hand, or the like) is extracted from the image, and the area is replaced with another image so as to be displayed on a display (for example, see patent document No. 1). A user interface is also known that receives, as an instruction for operating an application, the movement of a mouth or a hand captured by a video camera."}
{"text": "The present disclosure relates to coating metallic powder particles in general, and to processes for coating metallic powder particles with silicon in particular.\nAluminum alloy materials are widely used in many areas including the aerospace industry for their light weight structural properties and in heat transfer products due to high thermal conductivity. Additive manufacturing of aluminum alloy powders, in particular, via laser based additive processes, has drawn increasing attention. A significant challenge in the laser additive manufacturing of certain alloy powders (e.g., certain aluminum alloy powders) is the high degree to which energy in the form of light waves within a laser beam are reflected away from the alloy powder. As a result of the reflectance, the energy of the reflected light does not appreciably contribute to the additive process. In some instances, it may not be possible to sufficiently increase the intensity of existing laser equipment to overcome the reflectance issue. Even in those instances where the laser intensity can be increased, such a practice can create new detrimental issues; e.g., a higher intensity laser can overheat the powder and cause the powder particles to “ball up” and consequently create a non-uniform deposition layer. Also as a result of the reflectance, present additive manufacturing processes are sometimes limited to certain types of materials. With respect to aluminum alloys, for example, additive manufacturing processes are today typically limited to cast aluminum alloy compositions. These cast aluminum alloy compositions have low reflectance properties that allow them to be used in additive processes, but possess undesirable mechanical and thermal properties.\nIn addition, alloy powders adsorbed with water moisture can cause significant quality issues of the deposits made by additive manufacturing processes. These issues include porosity, cracks, and blisters.\nIt would be beneficial to resolve the aforesaid issues and thereby improve additive manufacturing processes and make it possible to additively manufacture certain alloy powders, including additional types of aluminum alloy compositions."}
{"text": "1. Field of the Invention\nThe present invention relates to integrated circuit test apparatus, more particularly, to test systems for so-called system LSIs in which a circuit for implementing one function is formed into a macro in a single integrated circuit together with another circuit.\n2. Description of the Related Art\nConventionally, circuits of, e.g., a CPU, a DSP, a DRAM, etc., are respectively mounted as independent devices on different boards. In recent years, as such circuits continue to shrink in feature size and increase in density, so-called system LSIs are increasing in which those circuits are formed as macros and integrated as verified function macros in a single integrated circuit together with other circuits. As an increasing trend, one integrated circuit includes two or more function macros.\nIn many cases, the layout of a function macro which assembles an independent device into a macro is fixed. It is, therefore, difficult to add a test support circuit such as a so-called logic scan circuit into a function macro later. An operation test for a system LSI including a function macro is often done by a function test of externally supplying an instruction or necessary signal, actually running the function macro, and extracting the processing result outside.\nFIG. 1 is a circuit diagram showing a conventional fundamental construction for running a function test on a function macro. FIG. 1 shows only the input side of the function macro as a representative. As shown in FIG. 1, a function macro 100 comprises input terminals IN0, IN1, . . . , INn−1, and INn. Selectors S0, S1, . . . , Sn−1, and Sn are so provided as to correspond to the respective input terminals IN0, IN1, . . . , INn−1, and INn.\nEach of the selectors S0 to Sn receives two inputs: a signal (instruction or necessary signal for a function test) input from the external input terminal of an integrated circuit including the function macro 100, in function test operation, and a signal input from a circuit of the integrated circuit other than the function macro 100 in normal operation of the integrated circuit. Each selector selects either of the signals, and supplies the selected one to the corresponding one of the input terminals IN0 to INn of the function macro 100. At this time, the selectors S0 to Sn switch the selection states of two inputs in accordance with control signals input to their control terminals.\nMore specifically, when the control signal represents that a test mode is set, each of the selectors S0 to Sn selects a signal input from the external input terminal of the integrated circuit in function test operation, and supplies the selected signal to the corresponding one of the input terminals IN0 to INn of the function macro 100. When the test mode is not set, i.e., a normal operation mode is set, each of the selectors S0 to Sn selects a signal input from a circuit of the integrated circuit other than the function macro 100, and supplies the selected signal to the corresponding one of the input terminals IN0 to INn of the function macro 100.\nIn the construction shown in FIG. 1, the terminal definition is changed by the selectors S0 to Sn to connect all the input terminals IN0 to INn of the function macro 100 to the external input terminals of the integrated circuit in test operation so as similarly to run an existing function test for an independent device on the verified function macro 100 integrated in the integrated circuit.\nFIG. 2 is a circuit diagram showing a conventional construction when one integrated circuit includes two different function macros. In the construction shown in FIG. 2, one integrated circuit includes two function macros A and B (101 and 102) and another user logic 103, to which necessary signals are supplied from external input terminals 107 of the integrated circuit to run function tests sequentially. The processing results in the function macros 101 and 102 are output from external output terminals 108 of the integrated circuit to obtain the results of the function tests.\nFor example, function tests are sequentially run on the function macros 101 and 102, and then a function test is run on the user logic 103. In running a function test on the function macro 101, a select signal is supplied from a function macro A select terminal 104 to the function macro 101. In running a function test on the function macro 102, a select signal is supplied from a function macro B select terminal 105 to the function macro 102. In running a function test on the user logic 103, a signal representing a test mode is supplied from a control signal input terminal 106 to the user logic 103.\nIn running a function test on the function macro 101, signals input from the external input terminals 107 of the integrated circuit are supplied to the function macro 101. At this time, signals input from some of the external input terminals 107 are supplied to the function macro 101 via selectors 109 each for selectively outputting either of the input signal and a signal processed by the user logic 103. The processing results in the function macro 101 are output outside from the external output terminals 108 of the integrated circuit via selectors 111 each for selectively outputting any of a signal processed by the function macro 101, a signal processed by the function macro 102, and an output signal in normal operation.\nIn running a function test on the function macro 102, signals input from the external input terminals 107 of the integrated circuit are supplied to the function macro 102. At this time, signals input from some of the external input terminals 107 are supplied to the function macro 102 via selectors 110 each for selectively outputting either of the input signal and a signal processed by the user logic 103. The processing results in the function macro 102 are output outside from the external output terminals 108 of the integrated circuit via the selectors 111. Note that the external output terminals 108 include a dedicated terminal for outputting only a signal processed by, e.g., the function macro 102.\nIn running a function test on the user logic 103, signals input from the external input terminals 107 of the integrated circuit are supplied to the user logic 103. Signals processed by the user logic 103 are supplied to the function macros 101 and 102 via the selectors 109 and 11 zero the processing results in the function macros 101 and 102 are output outside from the external output terminals 108 of the integrated circuit with or without the mediacy of the selectors 111.\nIn this construction, external terminals of the integrated circuit connected to only the function macros 101 and 102 are test terminals.\nAs shown in FIG. 1, an integrated circuit such as a conventional system LSI must be equipped with test terminals (external input terminals) for externally supplying various test signals to the input terminals of a function macro in order to run an existing function test even on a function macro integrated in the integrated circuit. This also applies to output terminals for outputting processing results in a function macro. For this purpose, a large number of external test terminals not used in normal operation must be provided on an integrated circuit, resulting in a highly costly integrated circuit.\nAs shown in FIG. 2, in an integrated circuit including function macros, the input/output terminals of function macros to be tested are selectively connected to the external test terminals of the integrated circuit to run function tests on the function macros sequentially. This prolongs the time required to test all the function macros.\nAs a means for preventing an increase in test time, as shown in FIG. 3, function tests are simultaneously run on two function macros 101 and 102 by simultaneously supplying a select signal for selecting function macros from a macro select terminal 112 to the function macros 101 and 102.\nIn this construction, however, all the input/output terminals of the function macros 101 and 102 must be connected to the external test terminals of the integrated circuit in test operation, to say the least. Thus the number of external test terminals of the integrated circuit must be increased. Such an increase in the number of test terminals greatly increases the cost."}
{"text": "1. Field of the Invention\nThe present invention relates to a golf ball with a cover wherein hit feeling and control property of the golf ball are improved.\n2. Description of the Related Art\nAs a cover material of a golf ball, an ionomer resin has widely been used (e.g. Japanese Patent Kokai No. 49-49727, etc.), heretofore. Particularly, in almost all of two-piece golf balls using a solid core, the ionomer resin is used as a cover material.\nThe reason for this is that the ionomer resin is superior in durability, cut resistance and impact resilience and is easily processed. Furthermore, the Ionomer resin is cheap in comparison with other cover materials.\nHowever, since the ionomer resin has very high hardness and stiffness, it is inferior in hit feeling and control property (easy spinning) in comparison with balata (transpolyisoprene), which is used as a cover base of a thread wound golf ball.\nAs a result, attempts to improving hit feeling and control property have been made by softening the ionomer resin with various means, however, satisfactory results are not obtained at present."}
{"text": "1. Field of the Invention\nThe present invention relates to an image pickup device suitable for a digital single lens reflex camera.\n2. Description of the Related Art\nIn a known single lens reflex camera, a light emission of a strobe is automatically controlled by a TTL system. In case of a single lens reflex camera using a silver-halide film, diffused light reflected from a photographic emulsion surface of the film is received and detected by TTL light receiving elements. In recent years, digital single lens reflex cameras in which an image pickup device such as CCD is used instead of a photographic film have been widely used. In a known image pickup device used in a digital camera, a micro lens is provided in front of a photoelectric transducer to effectively transmit object light incident through a photographing lens to the photoelectric transducer. To enhance the dynamic range, it has been proposed to provide a light reducing filter made of a material having a reflectance not more than 50% between a part of photoelectric transducers and a color filter (see Japanese Unexamined Patent Publication No. 2003-273338). Moreover, it has also been proposed to coat front and rear surfaces of a ferroelectric film in the form of a micro lens with transparent electrode layers so as to vary the reflectance (Japanese Unexamined Patent Publication No. 2001-332711).\nHowever, since the reflectance at the incident side surface of the image pickup device provided in a digital single lens reflex camera is small, a flash of light emitted from a strobe of a TTL auto strobe system and reflected from an object to be photographed is hardly reflected by the incident surface of the image pickup device and, hence, a small amount of the object light is made incident upon the TTL light receiving elements. Consequently, the output of the TTL light receiving elements is extremely small. As a result, the peak value in the photometering sensitivity distribution of the TTL light receiving elements is deviated from the center, and the photometering sensitivity distribution exhibits a steep curve and the output largely varies in accordance with the movement of the position of an exit pupil of a photographing lens. Also, the Av-response (proportional relationship between the actual aperture opening degree and the output) is not good. In the arrangement disclosed in Japanese Unexamined Patent Publication No. 2003-273338 in which the light reducing film is provided on a part of the photoelectric transducers, the output of the TTL light receiving elements largely varies depending on the position or shape of the object on the image surface, for the same object distance. In the solution disclosed in Japanese Unexamined Patent Publication No. 2001-332711 in which the front and rear surfaces of the micro lens are coated with the transparent electrode layers, a part of light to be made incident upon a photoelectric transducer is incident upon another photoelectric transducer, resulting in a deteriorated image quality."}
{"text": "1. Field of the Invention\nThis invention relates to a semiconductor integrated circuit device, for example, to a CMOS (Complementary Metal Oxide Semiconductor).\n2. Brief Description of the Prior Art\nConventionally, in well known semiconductor integrated circuit devices, for example, ones having a CMOS structure, latchup, which is characteristic of such CMOS device, occurs easily. That is, in adjacent elements (e.g., P channel and N channel MOS transistors), a PNPN or NPNP structured parasitic thyristor occurs and acts to induce malfunction by trigger action of some external factor.\nTo prevent such latchup, diffusion layers, referred to in general as guardrings, are provided in predetermined areas.\nWith regard to FIGS. 7 and 8 there is shown an N type well 2 in a P type silicon substrate 1, with P channel MOS transistors 20a and 20b having a source region (a P+ type diffusion region) 9 and drain regions (P+ type diffusion regions) 10, both provided in this N type well 2, (the two drain regions 10 being provided here with the common source region 9).\nOn the other hand, N channel MOS transistors 20c and 20d include a source region (an N+ type diffusion region) 4 and drain regions (P+ type diffusion regions) 5, both provided in the P type silicon substrate 1, (the two drain regions 5 being provided with the common source region 4). Furthermore, as seen in FIG. 7, a CMOS device (e.g., an inverter circuit) comprises two gate electrodes 7 commonly provided to the P channel MOS transistors 20a and 20b and to the N channel MOS transistors 20c and 20d respectively.\nA guardring region (an N+ type diffusion region) 8 is formed in a predetermined area so as to surround the P channel MOS transistors 20a and 20b. However, as seen in FIG. 7, this guardring region 8 is not provided in the area where the gate electrodes 7 cross over the guardring region 8 (i.e., the guardring region 8 is not continuously formed). Referring to reference numerals shown in FIG. 8, 3 is a field oxide film, 6 is a gate oxide film, 8a is an island shaped N+ type diffusion region (a part of the guardring region), 11 is a contact hole, and 12 is an electrode (or contacts) to maintain the guardring region at a supply voltage V.sub.cc.\nAccording to the above structure, as shown in a chart of FIG. 9, the conventional CMOS type 1M dynamic RAM has a higher resistance to latchup when the voltage is applied to the guardring (comparison 2) than when it is not applied thereto (comparison 1). However, the reliability is still insufficient, and the below noted problems remain.\n(1) In the area where the gate electrodes 7 cross over the guardring region 8 (directly below the gate electrodes 7), the guardring region 8 is mutilated and not continuously formed, so there is a possibility that latchup may occur through the area without the guardring.\n(2) Because there is the island shaped guardring region 8a as described above (one region in this case), a contact 12a must be defined to separately apply the supply voltage V.sub.cc also to this guardring region 8a (a certain area is required to provide this contact). It is difficult to make the contact in such a place with a minimum area so that taking excess area cannot be avoided."}
{"text": "The communications industry is rapidly changing to adjust to emerging technologies and ever increasing customer demand. This customer demand for new applications and increased performance of existing applications is driving communications network and system providers to employ networks and systems having greater speed and capacity (e.g., greater bandwidth). In trying to achieve these goals, a common approach taken by many communications providers is to use packet switching technology. Increasingly, public and private communications networks are being built and expanded using various packet technologies, such as Internet Protocol (IP).\nA network device, such as a switch or router, typically receives, processes, and forwards or discards a packet based on one or more criteria, including the type of protocol used by the packet, addresses of the packet (e.g., source, destination, group), and type or quality of service requested. Additionally, one or more security operations are typically performed on each packet. But before these operations can be performed, a packet classification operation must typically be performed on the packet.\nPacket classification as required for, inter alia, access control lists (ACLs) and forwarding decisions is a demanding part of switch and router design. The packet classification of a received packet is increasingly becoming more difficult due to ever increasing packet rates and number of packet classifications. For example, ACLs require matching packets on a subset of fields of the packet flow label, with the semantics of a sequential search through the ACL rules. IP forwarding requires a longest prefix match.\nKnown approaches of packet classification include using custom application-specific integrated circuits (ASICs), custom circuitry, software or firmware controlled processors, binary and ternary content-addressable memories (CAMs). The use of programmable software or firmware have advantages as they provide some level of flexibility, which becomes especially important as new protocols and services are added to existing network. Customer typically desire to use their existing hardware (e.g., routers, switches etc.) to support these new protocols and services. However, known software and firmware implementations are relatively slow, and typically place a performance bound which may be incompatible with new requirements. CAMs are increasingly being used in packet classification especially because of their performance. However, these content-addressable memories are typically expensive in terms of power consumption and space, and are limited in the size of an input word (e.g., 72, 144, etc.) on which a lookup operation is performed.\nVarious applications that use packet classification, such as Security Access Control, Quality of Service etc., typically need to perform many matches on source and destination port numbers, protocol and other header fields, etc.\nOne known approach for providing lookup capability on multiple data items is to use a separate CAM for each type of information being matched. For example, one CAM might be programmed with address information, another with quality of service information, and yet another with port information. Such an approach is expensive as it uses multiple CAMs. Needed are new methods and apparatus for implementing and using multiple virtual portions of a single associative memory, including a less costly approach for providing multiple CAM lookups on a data item.\nA known approach for accommodating a lookup word having a greater length than that allowed by a single lookup word of a CAM uses a single CAM and breaks the desired lookup string into different segments and performs multiple lookups. A CAM is programmed with parent entries (i.e., those for matching a first segment) and child entries (e.g., those for matching the remaining segments), wherein each child entry includes the address of its parent as part of its lookup string to ensure that only children entries of a matched parent entry are matched. Additionally, child and parent entries contain a different two bit tag to ensure a child entry is not matched when a parent entry is desired, and vice versa. A first segment of the string is provided as input to the CAM to produce a matching address (e.g., page and word address). A lookup word including a second segment of the string along with the matching address and child tagging is then provided as input to the CAM to produce a second matching address. This process can be repeated wherein the child entry becomes the parent entry for another child entry. This approach uses a single CAM, however, the use of the address of a parent entry as the prefix for a lookup of a child entry is inflexible, and typically requires a fixed number of bits which may be excessive.\nAnother known approach for accommodating lookup word having a greater length than that allowed by a single lookup word of a CAM uses tags each succeeding stored entry in a long comparand with its sequence number. The first entry is tagged “1”, the second entry tagged “2”, the third entry is tagged “3”, and so forth. After the first match is found, the second portion of the unknown is entered into the comparand register along with the tag for “2”, and if a match is found, the Status register is examined to see if the match address is equal to the previous match address +1. If it is, then the third portion of the unknown can be loaded into the comparand with the tag for “3”, and so forth. If the match address is not an increment from the previous address, that means the mach has failed, and the first matching location should be set to “SKIP” and a CMP V issued to find the next higher match."}
{"text": "This invention relates to treating the surface of metal sheets with electrolytes and the products thereby obtained. The resulting metal sheets have improved corrosion resistance and are suitable, among other uses, for lithographic applications. When used as supports in lithography, particularly if aluminum or its alloys are selected, such sheets exhibit improved adhesion for light sensitive coatings, improved printing plate run length, lessened wear by a press both in image and non-image areas, greater shelf life, and improved hydrophilicity in non-image areas.\nAnodization is an electrolytic process in which the metal workpiece is made the anode in a suitable electolyte. When electric current is passed, from a cathode through the electrolyte to the metal workpiece, the surface of the metal is converted to a form of its oxide having decorative, protective or other properties. The cathode is either a metal or graphite, at which the only important reaction is hydrogen evolution. This coating progresses from the solution side, outward from the metal, so the last-formed oxide is adjacent to the metal. The oxygen required originates from the electrolyte used.\nAnodic oxide coatings on aluminum may be of two main types. One is the so-called barrier layer which forms when the anodizing electrolyte has little capacity for dissolving the oxide. These coatings are essentially nonporous; their thickness is limited to about 13 .ANG./volt applied. Once this limiting thickness is reached, it is an effective barrier to further ionic or electron flow. The current drops to a low leakage value and oxide formation stops. Boric acid and tartaric acid are used as electrolytes for this process.\nWhen the electrolyte has appreciable solvent action on the oxide, the barrier layer does not reach its limiting thickness: current continues to flow, resulting in a \"porous\" oxide structure. Porous coatings may be quite thick: up to several tens of micrometers, but a thin barrier oxide layer always remains at the metal-oxide interface.\nElectron microscopic studies show the presence of small, close-packed cells of amorphous oxide throughout the oxide layer, generally peroendicular to the metal-oxide interface.\nSulfuric acid is the most widely used electrolyte, with phosphoric also popular. Anodic films of aluminum oxide are harder than air-oxidized surface layers.\nAnodizing for decorative, protective and adhesive bonding properties has used strong electrolytes such as sulfuric acid and phosphoric acid. U.S. Pat. No. 2,703,781 employs a mixture of these two electrolytes.\nU.S. Pat. No. 3,227,639 uses a mixture of sulfophthalic and sulfuric acids to produce protective and decorative anodic coatings on aluminum. Other aromatic sulfonic acids are used with sulfuric acid in U.S. Pat. No. 3,804,731.\nAs a post-treatment after anodization, the porous surface is sealed according to numerous processes to determine the final properties of the coating. Pure water at high temperature may be used. It is believed that some oxide is dissolved and reprecipitated as a voluminous hydroxide (or hydrated oxide) inside the pores. Other aqueous sealants contain metal salts whose oxides may be coprecipitated with the aluminum oxide.\nU.S. Pat. No. 3,900,370 employs a sealant composition of calcium ions, a water-soluble phosphonic acid which complexes with a divalent metal to protect anodized aluminum or anodized aluminum alloys against corrosion. Polyacrylamide has been proposed as a sealant.\nU.S. Pat. No. 3,915,811 adds an organic acid (acetic acid, hydroxy acetic acid, or amino acetic acid) to a mixture of sulfuric and phosphoric acids to form the electrolyte in preparation for electroplating the so-formed anodic aluminum coating.\nU.S. Pat. No. 4,115,211 anodizes aluminum by A.C. or superimposed A.C. and D.C. wherein the electrolyte solution contains a water-soluble acid and a water-soluble salt of a heavy metal. The water-soluble acid may be oxalic, tartaric, citric, malonic, sulfuric, phosphoric, sulfamic or boric.\nU.S. Pat. No. 3,988,217 employs an electrolyte containing quaternary ammonium salts, or aliphatic amines and a water-soluble thermosetting resin to anodize aluminum for protective, ornamental or corrosion resistant applications.\nThe advantages of anodized aluminum as a carrier for lithographic printing plates were early recognized. Processes employing as electrolytes sulfuric acid, phosphoric acid, mixtures of these, or either of these in succession have been proposed. Prior to anodizing the sheet may be roughened mechanically, electrochemically or chemically. The need for a subcoating prior to the application of a photosensitive layer was recognized to impart adhesion to the coating and hydrophilicity to the non-image areas. U.S. Pat. No. 3,181,461 uses an aqueous alkaline silicate treatment following the anodization step.\nU.S. Pat. No. 2,594,289 teaches (Col. 1, lines 42-54) that porous anodic films but not nonporous anodic films are suitable for lithographic purposes, \"since the porous film confers a better water receptive surface to the non-image areas of the plate and allows image-forming material to anchor effectively to the surface by penetrating the pores.\"\nU.S. Pat. No. 3,511,661, since disclaimed, describes aluminum sheet for a lithographic printing surface anodized in aqueous phosphoric acid having an anodic film with a cellular pattern of aluminum oxide having cells with porous openings of about 200 .ANG. to 700 .ANG. in average diameter and a surface with 10 to 200 mg per square meter of aluminum phosphate.\nU.S. Pat. No. 3,658,662 describes the electrochemical silication of a cleaned, etched aluminum plate to achieve a measure of hydrophilization.\nIn U.S. Pat. No. 3,902,976 a conventionally anodized aluminum sheet is electrolytically post-treated in an aqueous solution of sodium silicate to form a hydrophilic abrasion-resistant and corrosion-resistant layer suitable as a support for a presensitized lithographic sheet.\nU.S. Pat. No. 4,022,670 carries out anodization of aluminum sheets in an aqueous solution of a mixture of polybasic mineral acid such as sulfuric and a higher concentration of a polybasic aromatic sulfonic acid such as sulfophthalic acid to produce a porous anodic oxide surface to which a photosensitive layer may be directly applied.\nThere is described in U.S. Pat. No. 4,090,880, a two-step process whereby a cleaned aluminum sheet is first coated with an interlayer material such as alkali silicate. Group IV-B metal fluorides, polyacrylic acid, or alkali zirconium fluoride and then anodized conventionally in aqueous sulfuric acid. Enhanced shelf life when overcoated with diazo sensitizers is claimed.\nU.S. Pat. No. 4,153,461 employs a post-treatment with aqueous polyvinyl phosphonic acid at temperatures from 40.degree. to 95.degree. C. after conventional anodizing to a thickness of at least 0.2 um. The treatment provides good adhesion of a subsequently applied light sensitive layer, good shelf life and good hydrophilization of non-image areas after exposure and development as well as long press runs.\nU.S. Pat. Nos. 4,383,897; 4,399,021; 4,448,647; and 4,452,674 which are incorporated herein by reference also teach related electrolysis techniques.\nPlates of the above construction, particularly when the light sensitive layer is a diazo compound have enjoyed considerable commercial success. Nevertheless, certain improvements would be desirable. These include freedom from occasional coating voids, occasional unpredictible premature image failure on the press, faster, more dependable roll-up on the press and freedom from other inconsistencies. Still greater press life is desirable as well as a process that would be more economical than conventional anodizing followed by a second operation of sealing or post-treating in preparation for coating with a light sensitive layer.\nIn the case of protective and decorative applications, improved corrosion resistance and production economy over known anodizing processes is desired."}
{"text": "The present invention relates generally to an open superconductive magnet used to generate a uniform magnetic field as part of a magnetic resonance imaging (MRI) diagnostic system, and more particularly to such a magnet having a support structure for a high magnetic field strength.\nMRI systems employing superconductive or other type magnets are used in various fields such as medical diagnostics. Open magnets employ two spaced-apart superconductive coil assemblies (or other magnet pole assemblies) with the space between the assemblies allowing for access by medical personnel for surgery or other medical procedures during MRI imaging. The patient may be positioned in that space or also in the bore of annular coil assemblies. The open space helps the patient overcome any feelings of claustrophobia that may be experienced in a closed magnet design. Known open magnet support systems include \"C\"-shaped supports with the superconductive coil assemblies (or other magnet pole assemblies) being attached to the open top and bottom ends of the \"C\". Other known support systems include structural posts whose ends are attached to the two superconductive coil assemblies (or other magnet pole assemblies) including attachment to the coil forms (which hold the superconductive coils) of such superconductive coil assemblies. These support systems may be adequate for low field strength magnets, but a more robust support system is desired when the open magnets have a high field strength, including one Tesla or higher, where the attraction forces between the toroidal-shaped superconductive coil assemblies can be in excess of 160,000 pounds."}
{"text": "The present invention relates to a method for fabricating a composite laminate part, the parts obtained using this method, and the use of these parts in the automotive field to produce vehicle bodies, in the fields of domestic electrical appliances, industry in general and the building trade.\nThe new legislation enacted to reduce carbon dioxide emissions are forcing motor vehicle manufacturers to reduce the weight of motor vehicles in order to lower their fuel consumption."}
{"text": "This application is a National Stage of International Application No. PCT/EP2007/005557, filed Jun. 23, 2007, which claims priority under 35 U.S.C. §119 to German Patent Application No. 10 2006 038 124.6, filed Aug. 14, 2006, the entire disclosure of which is herein expressly incorporated by reference.\nThe present invention relates to a restraint system for a motor vehicle.\nIt is well known for motor vehicles to protect occupants from injuries by inflating an air bag in a short period of time, to catch the occupant who is displaced in a forward direction. For this purpose, it is important that the air bag quickly reaches its effective volume, and it is therefore necessary to produce relative large amounts of gas (e.g., in a pyrotechnical manner) and to introduce them into the gas bag. Recently, for dimensioning the gas volume to be added, occupant/and or vehicle parameters have been evaluated, so as not to unfold the gas bag to its full size for example in the case of a so-called “out of position” of the occupant. However, the necessity for adjustment to all possible load conditions and the relatively cost-intensive control and regulation technique have made this feature problematic.\nIntensified efforts have been made, therefore to develop so-called self-regulating systems which can adjust automatically to the corresponding load conditions. For example, venting apertures have been provided, which close automatically when the inner pressure of the gas bag becomes too large.\nGerman patent document DE 2 302 737 discloses a restraint system comprising a two-layer gas bag, where the gas is guided only between the two layers, so that a complete gas cushion does not result; rather, a spherical annular supporting structure. The unfolding of the support structure thereby occurs in the transverse direction of the abutting gas cushions. On the other hand, European patent document EP 0589 059 B1 shows furthermore, that it is necessary during the unfolding of the two-layer gas bag, to suck ambient air into the interior so as to overcome the negative pressure.\nBoth of these systems have in common the feature that less gas volume is necessary to unfold the gas bag to its full size, due to the gas bag's being formed with two layers. The temperature and the pollution can thereby be reduced.\nFinally, UK patent document GB 1 420 226 A discloses a restraint system for a motor vehicle, where a tubular supporting structure is provided in the interior of the two-layer gas cushion, the longitudinal extension of which exceeds the dimensions of the transverse extension in the inflated (that is, active) state. The supporting structure unfolds due to its particular geometry mainly in the direction of its longitudinal extension.\nAn automatic adjustment of the gas bag dependent on the respective load condition (e.g., dependent on the respective occupant or his or her position) is not described.\nIt is therefore an object of the present invention to improve the generic restraint system in such a manner that its size adjusts to the respective load condition.\nThis and other objects and advantages are achieved by the restraint system according to the invention, which when it impacts upon an obstacle during unfolding, does not achieve the stability it would have when fully unfolded and/or the final volume or the final extension, due to the particular geometric form of the supporting structure (that is, due to its longitudinal extension exceeding its cross section). The supporting structure has a different stability during unfolding, due to its longitudinal formation. While the supporting structure is initially rather unstable (that is, can be easily hindered during unfolding), the completely unfolded supporting structure achieves full stability (that is restraining force). This means that, if the supporting structure impacts upon an obstacle during the unfolding, as is the case for example with occupants leaning forward (out of position), unfolding can be stopped more easily or deflected due to the supporting structure which is still unstable. These small forces effect a low pressurization of the occupant.\nAs used herein, the term “supporting structure” refers to a structure which is similar to a skeleton (in contrast to a conventional spherical gas bubble—a gas bag), and which comprises a restraint effect in the fully unfolded state comparable to the conventional gas bag. The provision of a substantially more complex inflatable supporting structure according to the invention, for example a branched tree structure not only reduces the necessary gas volume, but also diminishes the force peaks acting on the occupant during unfolding, if he is effectively “in the way” of the unfolding supporting structure. In contrast, the gas amount in the conventional gas bubble is on a substantially higher level from the start of the activation until the full unfolding, so that obstacles in the unfolding path are put under more pressure independent of the unfolding state.\nThe supporting structure can preferably be filled with gas, and it is unimportant whether it is filled with gas generated in a pyrotechnical manner or with gas from a pressure vessel or the ambient air. The ambient air can be guided into the carrier volume enclosed by the supporting structure by suction through the negative pressure occurring during the unfolding. The gas volume can be reduced further in this manner.\nIf several supporting structures are provided, which are connected to one another or to the environment in a fluidic manner, a restraint system can be formed, which is constructed in a net-like, tree-like or supporting frame-like manner. Thus, supporting structures proceeding in the longitudinal direction of the vehicle can be connected to supporting structures proceeding transversely to the transverse direction of the vehicle. The unfolding of the transverse supporting structures thereby nevertheless takes place in the direction of the longitudinal extension of the respective supporting structure.\nThe supporting structures can be connected by flexible sheets, in particular flexible textile sheets. In this case, the flexible sheet can be fixed between two adjacent supporting structures.\nIf several supporting structures are connected to one another by a flexible sheet, the gas can be used for restraining action in the resulting enclosed carrier volume enclosed. The gas in the carrier volume serving to restrain occupants can be heated or supplemented by means of a heating device (e.g., an ignition tablet or a small gas generator step, which develops heat with a possibly very low gas volume). As a result, the gas volume or the inner pressure, and thus the restraint action, increases therewith correspondingly.\nIf a supporting structure impacts upon an obstacle during unfolding, the gas flowing into the supporting structure can be distributed to other adjacent supporting structures or into the environment, by gas redistribution components. The inner pressure generated by the flow in the supporting structure can also be reduced by increasing its cross section. This can take place in such a manner that tear seams break down at a certain pressure, so that the supporting structure can increase in its transverse direction. A pressure-relief valve can also be provided, which opens when an interior pressure is reached which is too high.\nWith conventional compact gas bags, it was previously necessary to refrain from directing venting apertures towards the occupant, because gas temperatures are reached which are too high. However, with the system according to the invention, as cold ambient air flows within the carrier volume, and not a pyrotechnically produced gas, and venting apertures can be directed towards the occupant by the reduced temperature. The air bag dampening can thus be adapted for various environmental conditions, by sealing a venting aperture cross section between occupant and air bag.\nThus, this contact surface and its sealing with persons having a higher volume (and usually a higher weight) is larger, so that a stronger restraint action is achieved hereby. With more severe accidents, the contact surface and thus the restraint action is also increased by a stronger immersion of the occupant into the air bag. This principle also permits a variable air bag dampening for belted and unbelted occupants, because in the unbelted state, the occupant will be immersed into the air bag with less force.\nIf the venting apertures are formed as perforations (that is, many small apertures), as for example with a textile net, a projection area corresponding to the measurements of the occupant can be closed. The reproducibility of the results increases on average with many small venting apertures.\nIn a preferred embodiment, at least one of the supporting structures extends at least partially in an overlap of a hard structure, as for example a supporting column. It is also possible to provide a structure which unfolds between the occupants as an interaction bag.\nSo as to be able to manufacture a complex supporting structure with its inflatable components and the flexible sheets, the one-piece woven technique is recommended. This technique distinguishes itself in that, on one and the same device, a double layer can be manufactured for the inflatable structures, and one layer for the flexible sheets, or a three-dimensional structure can be woven.\nOther objects, advantages and novel features of the present invention will become apparent from the following detailed description of the invention when considered in conjunction with the accompanying drawings."}
{"text": "This invention relates to a semiconductor memory device as a multi-port DRAM having split SAM registers.\nAn example of the conventional DRAM of this kind is shown in FIG. 1. As seen from this figure, the cell array constituting the RAM section is divided into a cell array where the value of the Most Significant Bit (MSB) value in the column address is \"1\" (hereinafter referred to as a high order cell array 1U) and a cell array where the value of the MSB is \"0\" (hereinafter referred to as a low order cell array 1L), and includes high order and low order SAM registers 2U and 2L in correspondence with the respective cell arrays 1U and 1L. Reference numeral 3U represents bit line pairs constituting data transfer paths between the high order cell array 1U and the high order SAM register 2U, and reference numeral 3L represents bit line pairs constituting data transfer paths between the low order cell array 1L and the low order SAM register 2L. In addition, transfer gates 4U and 4L are inserted into the bit line pairs 3U and the bit line pairs 3L, respectively.\nThe details of the relationship of the transfer gate 4U, the bit line pairs 3U and the high order SAM register 2U and the relationship of the transfer gate 4L, the line pairs 3L and the low order SAM register 2L are shown in FIGS. 2A and 2B, respectively.\nIn FIG. 2A, reference numeral 2UC represents cells constituting the high order SAM register 2U, and reference numerals 3U1 and 3U0 represent bit lines constituting respective bit line pairs 3U. The transfer gate 4U is comprised of a plurality of transistors 4UT. In addition, transistors 4UT are respectively inserted into bit lines 3U1 and 3U0 constituting respective bit pairs. Similarly, in FIG. 2B, reference numeral 2LC represents cells constituting the low order SAM register 2L, and reference numerals 3L1 and 3L0 represent bit lines constituting respective bit line pairs 3L. The transfer gate 4L is comprised of a plurality of transistors 4LT. In addition, transistors 4LT are respectively inserted into bit lines 3L1 and 3L0 constituting respective bit pairs.\nIn the above-mentioned configuration, by respectively allowing the transistors 4UT and 4LT to be turned ON and OFF, only the bit line pairs 3U on the high order side become active, resulting in the state where data transfer from the high order cell array 1U only to the high order SAM register 2U can be carried out. In contrast, by respectively allowing the transistors 4UT and 4LT to be turned OFF and ON, only the bit line pairs 3L on the low order side become active, resulting in the state where data transfer from the low order cell array 1L only the low order SAM register 2L can be carried out. Accordingly, by alternately placing both registers 2U and 2L on the high order and low order sides in a standby state and in an operational state to carry out the split transfer, continuous access from the SAM side can be provided.\nAs stated above, by employing a configuration in which the RAM section cell array is halved by the value of MSB of the column address and SAM registers are provided in correspondence with respective cell arrays, split transfer is performed every 1/2 row, permitting random access every 1/2 row from the SAM side. Thus, so-called window control can be efficiently carried out.\nNamely, FIG. 3 represents a memory space of the RAM section and FIG. 4 represents a display image subjected to mapping in correspondence with the memory space. In these figures, each rectangular region indicates one word, and the figure in the word indicates an address number It is to be noted that when a specific word is designated, a notation is employed in the following description such that alphabetical symbol \"AD\" is attached before the address number. For example, in the case of the word of the address number \"1\", that word will be called \"word AD1\".\nLet as now consider the case where a window is set in FIG. 4 in the region extending over the words AD6, AD7, AD10, AD11, AD14 and AD15 to control the window thus set.\nInitially, in the case of an access every row, in order to make an access to the region extending over the words AD6, AD7, AD10, AD11, AD14 and AD15, the region including the words AD5, AD8, AD9, AD12, AD13 and AD16 must be accessed, resulting in redundant bits of at least one byte, which are useless for the window control.\nOn the contrary, if an access every 1/2 row can be made, it is not required to make an access to the region including the words AD5, AD8, AD9, AD12, AD13 and AD16. It is sufficient only to access the region including the words AD6, AD7, AD10, AD11, AD14 and AD15. As a result, the efficiency is improved.\nHowever, in a conventional multi-port DRAM having split registers, there is the problem that when a frame buffer is constructed, mapping on a display image is restricted. Namely, in the case of FIG. 4, by interchangeably making accesses to the SAM registers 2U and 2L on the respective high order and low order sides, continuous access with respect to the words AD6, AD7, AD10, AD11, AD14 and AD15 can be provided.\nOn the other hand, FIG. 5 shows a display image subjected to mapping in another form. Let now consider the case where the control of a window set in the region extending over the words AD4, AD5, AD7, AD8, AD10 and AD11 is carried out.\nIn the case of carrying out a continuous split transfer, accesses to the SAM registers 2U and 2L on the respective high order and low order sides are necessarily made. However, in the case of mapping shown in FIG. 5, if such a split transfer is used, since continuous words AD5 and AD7 and continuous words AD8 and AD10 belong to the same cell array 1U or 1L, it is inevitable for maintaining continuity of those words to transfer words AD6 and AD9 to the registers 2U and 2L even if no readout operation is carried out from the registers 2U and 2L. For this reason, such a mapping as shown in FIG. 5 cannot be carried out.\nAs stated above, in the conventional multi-port DRAM, mapping is disadvantageously restricted in constituting a frame buffer."}
{"text": "A variety of valves are positioned within body vessels in animals to permit substantially unidirectional fluid flow through the body vessel from one body location to another. For example, native valves within the heart and veins function to regulate the direction of blood flow within the blood vessels of the body. Heart valves positioned within the heart direct the flow of blood to and from other organs and pump oxygenated blood to the rest of the body. Venous valves are typically bicuspid valves positioned at varying intervals within veins to permit substantially unidirectional blood to flow toward the heart. Body vessels such as veins transport blood to the heart and arteries carry blood away from the heart.\nDefects or injury to valves within a body vessel can compromise valve function, thereby disrupting the normal flow of fluid within the body vessel. For example, compromised valve function within a blood vessel may result in an undesirable amount of retrograde fluid flow within the blood vessel across a valve therein, and compromise the unidirectional flow of fluid across the valve. Retrograde fluid flow refers to fluid flow opposite the primary direction of fluid across the valve. For example, for a venous valve, retrograde fluid flow is blood flow away from the heart. Methods of treatment and medical devices suitable for implantation within a body vessel are provided herein.\nIn the condition of venous valve insufficiency, the venous valve leaflets do not function properly for a variety of reasons. For instance, the vein may become too large in relation to the leaflets so that the leaflets cannot come into adequate contact to prevent backflow (primary venous valve insufficiency), or clotting within the vein may thicken the valve leaflets (secondary venous valve insufficiency). Incompetent venous valves can result in symptoms such as swelling and varicose veins, causing great discomfort and pain to the patient. If left untreated, venous valve insufficiency can result in excessive retrograde venous blood flow through incompetent venous valves, which can cause venous stasis ulcers of the skin and subcutaneous tissue. Venous valve insufficiency can occur, for example, in the superficial venous system, such as the saphenous veins in the leg, or in the deep venous system, such as the femoral and popliteal veins extending along the back of the knee to the groin.\nVarious implantable medical devices are advantageously inserted within various portions of the body to treat conditions related to compromised valve function within a body vessel. Minimally invasive techniques and instruments for placement of intraluminal medical devices have been developed to treat and repair undesirable conditions within body vessels, including treatment of conditions that affect blood flow such as venous valve insufficiency. Various percutaneous methods of implanting medical devices within the body using intraluminal transcatheter delivery systems can be used to treat a variety of conditions. One or more intraluminal medical devices can be introduced to a point of treatment within a body vessel using a delivery catheter device passed through the vasculature communicating between a remote introductory location and the implantation site, and released from the delivery catheter device at the point of treatment within the body vessel. Intraluminal medical devices can be deployed in a body vessel at a point of treatment and the delivery device subsequently withdrawn from the vessel, while the medical device retained within the vessel to provide sustained improvement in valve function or to increase vessel patency. For example, published U.S. Patent Application US2004/0225352, filed Mar. 10, 2004 by Osborne et al. and incorporated herein by reference in its entirety, describes implantable medical devices comprising a valve for regulating fluid flow through a body vessel. The medical devices may include a valve leaflet attached to a radially-expandable support frame, and configured to permit both fluid flow in a first direction and a controlled amount of fluid flow in a second direction.\nOne challenge for development of an implantable prosthetic valve with the venous system is mitigating thrombus formation that can occlude the vessel and/or lead to loss of functionality of the valve structures that regulate blood flow. In contrast to the arterial system, the lower flow rates in the deep veins of the legs and feet can lead to stagnation of blood in the pockets about the bases of the leaflets or valve structure due to the inability of the blood to be flushed and refreshed thereabout. The pockets can fill with thrombus that compromises the ability the leaflets or valve structure to open and close in response to antegrade and retrograde flow (i.e., pressure differentials across the valve). For example, fibrinogen absorbed on to the surface of an implanted prosthetic valve can form a layer that triggers the biochemical pathway leading fibrin deposition, platelet aggregation, and thrombus formation.\nRemodelable materials, such as extracellular matrix (ECM) materials, can be used to provide a non-thrombogenic surface in an implantable prosthetic valve. Prosthetic valves desirably include valve leaflets formed from a remodelable material such that, upon implantation, the remodelable material can become vascularized to form a permanently non-thrombogenic leaflet surface. Small intestinal submucosa (SIS) is a commercially available ECM material (Cook Biotech Inc., West Lafayette, Ind.) derived from a porcine source and processed to retain remodelability. While the ability of valve leaflets made of ECM materials to remodel has been demonstrated clinically, the surface of the newly-implanted SIS can be vulnerable to thrombus formation, particularly in the pocket regions. Because remodeling is a process that can take 30 days or longer, depending on the environment, thrombogenicity has remained a clinical issue to be addressed when using remodelable biomaterials. Higher levels of both antegrade and retrograde fluid flow across the ECM material may enhance the remodeling process, for example by preventing or reducing stagnation of fluid in contact with the ECM material that may lead to thrombus formation. However, high levels of retrograde fluid flow that promote remodeling of the ECM material may reduce the clinical effectiveness of the valve design.\nWhat is needed are methods and devices for reducing undesirable levels of retrograde fluid flow across a valve within a body cavity while permitting desirably high levels of remodeling of an ECM material within the valve. Methods for providing a flow regulating medical device comprising an ECM material within a body vessel that permit both remodeling of the ECM material and a therapeutically effective level of retrograde fluid flow across the medical device are particularly desirable."}
{"text": "Polysilicon formed on a field oxide layer may be used as a resistor element for use in the semiconductor integrated circuit. In such a case, a plurality of polysilicon members having the same shape and the same characteristics may be formed as a resistor element unit for the purpose of providing a desired resistance value. The number of polysilicon members connected in series or in parallel is then adjusted to adjust the resistance value. The size and shape of such polysilicon members are required to be highly precise. The polysilicon members may thus need to be manufactured under the same process conditions. Namely, the conditions of photolithography and etching performed during the wafer process may need to be maintained constant with respect to each polysilicon member.\nWhen a plurality of polysilicon members having the same shape are arranged in matrix form at constant intervals, polysilicon members situated in the interior of the matrix are formed under the same process conditions because each of these polysilicon members has identical positional relationships with the surrounding adjacent polysilicon members. The polysilicon members situated at the periphery of the matrix, however, have adjacent polysilicon members only on one side thereof, and, thus, are subjected to different process conditions than the conditions applied to the interior polysilicon members which are surrounded by adjacent polysilicon members on the four sides thereof. Accordingly, the polysilicon members situated at the periphery are provided, as dummy members unused for actual circuit operations, only for the purpose of achieving constant process conditions for the interior polysilicon members.\nFIG. 1 is a plan view illustrating an example of a related-art resistor-element unit. FIG. 2 is a cross sectional view of the resistor-element unit illustrated in FIG. 1 as taken along a line A-A′.\nAs illustrated in FIG. 1 and FIG. 2, a resistor-element unit 10 includes a P-type substrate 13, a field oxide layer 14 formed on the P-type substrate 13, and a plurality of polysilicon members 11 and 12 formed on the field oxide layer 14. The polysilicon members 11 and 12 are arranged in matrix form at constant intervals which are an interval “a” in the vertical direction and an interval “b” in the horizontal direction. The polysilicon members 12 are dummy polysilicon members situated at the periphery of the matrix, and the polysilicon members 11 are resistor-element polysilicon members situated in the interior of the matrix.\nThe resistor-element polysilicon members 11 situated in the interior of the matrix are formed under the same process conditions because each of these polysilicon members has identical positional relationships with the surrounding adjacent polysilicon members. Accordingly, the resistor-element polysilicon members 11 have the same resistor-element characteristics. A desired resistor value can thus be obtained by connecting a desired number of resistor-element polysilicon members 11 in series or in parallel.\nThe dummy polysilicon members 12 situated at the periphery of the matrix, however, have adjacent polysilicon members only on one side thereof, and, thus, are subjected to different process conditions than the conditions applied to the interior resistor-element polysilicon members 11 which are surrounded by adjacent polysilicon members on the four sides thereof. Accordingly, the dummy polysilicon members 12 situated at the periphery are formed, as dummy members unused for actual circuit operations, only for the purpose of achieving constant process conditions for the interior polysilicon members.\nIn the resistor-element unit 10 formed as described above, the dummy polysilicon members 12 situated at the periphery are not used as part of an actual operating circuit, which means that there are needless elements and areas provided in the semiconductor integrated circuit chip. This gives rise to problems such as an increase in chip size and a cost increase."}
{"text": "Multiple application platforms are available in the computing industry today. Examples include AS/400, AIX, and z/OS operating systems (marketed by International Business Machines (IBM) Corporation of Armonk, N.Y.), and Windows 2000 (marketed by Microsoft of Redmond, Wash.). Another example of an application platform is IBM WebSphere Application Server, which is a middleware platform for distributed applications that is based on Java and can thus execute on multiple operating system platforms. Large scale computing solutions of today often consist of multiple computing applications, executing on different platforms, that are networked together as set of components that work together cooperatively to form a larger system or environment that addresses a particular business requirement.\nSince the requirements of operating systems and middle-ware application platforms typically differ, each system maintains its own user registry, which includes a list of users and associated information, such as user Ids and passwords, used to authenticate a user when access to the network is requested. A user may be a human user, or may be a software process assigned a local user identity, such as a print server. Each platform typically has its own administrative tools that allow a system administrator to add, delete or modify user identities in the user registry. With a heterogeneous network that has several different operating systems, this means that the system administrator must learn and become proficient in several different tools which handle identity management in their respective realms (e.g., platforms).\nIn view of the above, a need exists in the art for a novel approach to authenticated identity propagation and translation within a multi-component transaction processing environment to, for example, facilitate run-time inter-operation between components that are employing disparate security and user registry services, and to do so with performance characteristics and within programming model constraints that are suitable for use in the mainframe components of the environment"}
{"text": "This invention relates to a pneumatic device which is subject to positive feedback and, more particularly, to a sloped nozzle within the device designed for reducing the positive feedback.\nAs disclosed hereinafter, the present invention is particularly useful in conjunction with reversing pneumatic amplifiers which are designed to reverse the action of sensors with respect to the branch line output of the sensing system although the present invention may have applications wider than in the described reversing relay.\nAs described herein, the present invention is useful in connection with the Velocitrol CP980 manufactured by Honeywell Inc. and described in U.S. Pat. No. 4,182,486. The Velocitrol includes a velocity sensor which comprises a pair of nozzles mounted generally transversely to the air moving through a duct. One of the nozzles, the primary nozzle, is designed to be connected to a source of pressure and the other of the nozzles, the secondary nozzle, is designed to receive air issuing from the primary nozzle dependent upon the velocity of air moving through the duct. As the velocity of the air moving through the duct increases, the amount of air received by the secondary nozzle decreases and as the velocity of the air moving through the duct decreases the amount of air received by the secondary nozzle increases. This Velocitrol controller thus is a reverse acting device.\nA pneumatic device is necessary, therefore, to reverse the action of the velocity sensor to make it direct acting. However, it was found that the reversing pneumatic device selected for converting the Velocitrol to a direct acting controller was subject to positive feedback. This problem of positive feedback gave rise to the present invention for eliminating positive feedback in pneumatic devices."}
{"text": "1. Field of the Invention\nThe present invention relates to a device for the automatic harvesting of mushrooms.\n2. Description of the Related Art\nWhen mushrooms are harvested they are manually picked. The person that harvests the mushrooms must make his choice of the desired dimension by sight. The stem of the mushroom must then be cut off. This happens usually with a hand-held knife. The mushroom is then put down in a box or carton. This way of harvesting requires great skill and is very labor intensive. There have been trials with fully automatic harvesting machines, but these have proven to be failures. There are also machines that harvest the whole cultivating bed in one passage. This way of operation has the disadvantage that the harvest is a mixture of too large to too small. In this way the cultivating bed is not completely used."}
{"text": "The present disclosure relates to an air-turbine starter used to start gas turbine engines, and more particularly to an aerodynamic flowpath thereof.\nMany relatively large turbine engines, including turbofan engines, may use an air turbine starter (ATS) to initiate gas turbine engine rotation. The ATS is typically mounted on the accessory gearbox which, in turn, is mounted on the engine or airframe. Consequently, the ATS is installed in the aircraft at all times even though active operation may occur only for a minute or so at the beginning of each flight cycle, along with occasional operation during engine maintenance activities.\nThe ATS generally includes a turbine section coupled to an output section within a housing. The turbine section is coupled to a high pressure fluid source, such as compressed air, to drive the output section through a gear system. Thus, when the high pressure fluid source impinges upon the turbine section, the output section powers the gas turbine engine."}
{"text": "1. Field of the Invention\nThe present invention relates to a sample-and-hold circuit.\n2. Description of the Related Art\nA sample-and-hold circuit as shown in FIG. 3, which constitutes a part of an active matrix type liquid crystal driver circuit is known. Analog signals are input to the sample-and-hold circuit through signal lines V1 to Vn. The sample-and-hold circuit is provided with capacitors C1 to Cn and capacitors D1 to Dn for sampling and holding input voltages from the signal lines (the capacities of those capacitors include wire capacitances and stray capacitances). Ends on one side of the capacitors C1 to Cn are connected to the signal lines V1 to Vn through the analog switches S1 to Sn which are turned on and off by control signals SS1 to SSn, respectively, and the other ends of the capacitors C1 to Cn are all connected to ground. Ends on one side of the capacitors D1 to Dn are connected to the ends on one side of the capacitors C1 to Cn through analog switches H1 to Hn which are turned on and off by a control signal HS. The other ends of the capacitors D1 to Dn are connected to ground. Noninverting input terminals of operational amplifier circuits P1 to Pn are connected to the ends on one side of the capacitors D1 to Dn, respectively, while the inverting input terminals are connected to output terminals thereof, respectively. The output terminals of the operational amplifier circuits P1 to Pn are tied to a common line. Output control signals CN1 to CNn are input respectively to the operational amplifiers.\nWhen high-level control signals SS1 to SSn are input respectively to the analog switches S1 to Sn, the analog switches turn on, for this input period, so that the capacitors C1 to Cn are connected to the signal lines V1 to Vn and charged. When the control signals SS1 to SSn go to the low level to turn off the switches S1 to Sn, the signal voltages just before that are held in the capacitors C1 to Cn, respectively. Then, when a high-level control signal HS is input, the switches H1 to Hn turn on for the input period, so that the capacitors D1 to Dn are respectively connected in parallel with the capacitors C1 to Cn and charged. After the control signal HS goes to the low level to turn off the switches H1 to Hn, the capacitors D1 to Dn hold voltages corresponding to the voltages heretofore held by the capacitors C1 to Cn, respectively.\nWhen specified, output control signals CN1 to CNn are applied to the operational amplifier circuits P1 to Pn, respectively, their outputs are made active, and the voltages supplied to the noninverting input terminals are output, respectively. Therefore, by using control signals CN1 to CNn, it is possible to determine which of the voltages held by the capacitors D1 to Dn are output. For example, in case that a control signal CN1 is applied to the operational amplifier circuit P1, the output of the operational amplifier circuit P1 is made active, and the voltage held in the capacitor D1 is output as the output Vo. By changing over the control signals CN1 to CNn respectively, the voltages held in the capacitors D1 to Dn are output sequentially as the output Vo from the operational amplifier circuits P1 to Pn.\nHowever, in such a sample-and-hold circuit, operational amplifier circuits are provided as buffer circuits on a one-to-one correspondence with the capacitors D1 to Dn. As a result, the number of devices constituting the sample-and-hold circuit is large. Accordingly a large chip area is required when this circuit is formed as an integrated circuit."}
{"text": "The present invention pertains to a device for sorting coins of a set of coins with a coin feed device, with a sorting plate joining the coin feed device, having a guide edge, a guide strip extending along the guide edge as well as a series of sorting openings graduated for different coin diameters with cross sections increasing in the direction of conveying, a set of sorting plates for such a device, as well as to a process for converting such a device to different sets of coins.\nCoins are defined within the framework of the present invention as disk-shaped, mostly round objects which represent a monetary value. Thus, besides legal means of tender, the term coins also covers value chips, tokens, coin blanks, round blanks for coins, etc. A device for conveying coins on the sorting plate, e.g., a conveyor belt, is typically set up. The coins are conveyed between the conveyor belt and the sorting plate. The conveyor belt, including its guide elements, may be introduced detachably or pivotably, so that access to the sorting plate is possible. When a coin is being conveyed on a sorting plate, it moves along the guide edge. Each coin moving on the sorting plate is continuously supported by the guide strip. The sorting openings have different opening cross sections with increasing opening cross section q measured at right angles to the direction of conveying. A coin whose diameter is smaller than the sum of the cross sections q and the width d of the guide strip is supported only on the (comparatively very narrow) guide strip (whose width is, e.g., 0.05 to 0.4 times q) on one side and it consequently falls through the sorting opening. The sorting openings of a sorting plate are graduated with increasing opening cross sections in the direction of conveying. As a result, sorting of all coins of the set of coins is achieved. A means for counting the unsorted or sorted coins may be additionally provided within the framework of such a device. Furthermore, a device for sorting out counterfeit coins may be provided between the coin feed device, typically a rotary table, and the sorting plate. Reference is made in this connection only, e.g., to the literature reference EP 0 189 429 B1 which corresponds to U.S. Pat. No. # 4,681,204.\nDevices of the design described in the introduction have been known from, e.g., the literature reference DE 196 33 518 C1. The drawback of the devices known in this respect is that the device is delivered for a certain set of coins and a conversion of the device to another set of coins is possible with very great installation and adjustment efforts only. This is disturbing especially in cases of a changeover to another currency, e.g., the changeover from European national currencies to the euro.\nTherefore, compared with the state of the art, the basic technical object of the present invention is to propose a device for sorting coins which can be converted to another set of coins simply and with little effort.\nTo accomplish this technical object, the present invention teaches that the sorting plate is designed as an interchangeable sorting plate. A sorting plate that can be removed and mounted with simple means, preferably without any tool, is called an interchangeable sorting plate. It is achieved with the present invention that, e.g., the changeover from European national currencies to the euro can be carried out without an appreciable installation effort and consequently without an appreciable loss of time. Moreover, it is possible to change back from one set of coins to the other one and back again to the first one in an extremely short time during a transition period. Finally, the present invention makes possible the very simple and inexpensive replacement of used sorting plates.\nA preferred embodiment of the present invention is characterized in that the interchangeable sorting plate has at least one adjusting surface, wherein the interchangeable sorting plate can be inserted in the course of its insertion into the device by means of an adjusting element which interacts with the adjusting surface and is a rigid part of the machine. This embodiment of the present invention has a number of advantages concerning the problem of the accurate adjustment of the interchangeable sorting plate. On the one hand, the use of fits with narrow tolerances, which could lead to complications during installation, e.g., due to canting, etc., is avoided. On the other hand, the need for first adjusting and fixing an interchangeable sorting plate that may have been inserted loosely before putting the device into operation is eliminated.\nThe interchangeable sorting plate preferably has at least one fixing element, e.g., a thumbscrew, for fixing the interchangeable sorting plate inserted on a carrier plate. Any other fixing connections which can be easily handled manually, e.g., snap-in connections or quick-acting closures, are, of course, also possible instead of a thumbscrew.\nThe present invention may be specifically designed such that the adjusting surface is bent at an angle in relation to the longitudinal extension of the interchangeable sorting plate. It is achieved as a result that in the course of the insertion of the interchangeable sorting plate, the interchangeable sorting plate is brought automatically into the desired position by the interaction between the adjusting surface and the adjusting element. Correspondingly, the orientation of the adjusting surface shall be selected according to the desired end position. It is recommended for this that the end of the adjusting surface have a partial rest surface when viewed in the direction of installation. The accurate positioning is then carried out in practice in the simplest case such that an operator pushes in the interchangeable sorting plate and the adjusting surface is pushed, suitably guided manually, along the adjusting element into the end position of the interchangeable sorting plate. Additional adjusting surfaces may be arranged in parallel or in opposite directions (mirror symmetrically to a plane extending at right angles to the longitudinal extension of the interchangeable sorting plate) along the longitudinal extension of the interchangeable sorting plate. A pair of adjusting surfaces relieves an operator of the need to push along an adjusting surface against the adjusting element. The movement into the desired end position takes place automatically. The adjusting surface may be arranged in a closed recess of the interchangeable sorting plate or in a recess of the interchangeable sorting plate that is open (e.g., on one side). The adjusting element, which is a rigid part of the machine, may now be designed as a pin engaging the recess. It is recommended in this case to provide a xe2x80x9cclearancexe2x80x9d in the area of the end position or partial rest surface to avoid problems because of abraded particles or other contaminants, which may become deposited in the recess. xe2x80x9cClearancexe2x80x9d means that the free cross section within the recess is larger in all areas than the cross section of the adjusting element and thus it permits the adjusting surface or the partial rest surface to move off from the adjusting element in all positions of the interchangeable sorting plate. In the case of two recesses with adjusting surfaces which are arranged in each recess on one side and extend in opposite directions, accurate positioning can be achieved even despite the provision of xe2x80x9cclearancexe2x80x9d in every individual recess without the need for a special manual guiding during the installation. Only the adjusting elements are to be inserted at spaced locations such that moving off of the adjusting elements from the adjusting surfaces or partial rest surfaces is no longer possible in the inserted state.\nAn especially simple fixation of the interchangeable sorting plate on the carrier plate, which can be carried out especially rapidly, is achieved if at least one laterally projecting holding strap, which can be introduced into a holding recess that is a rigid part of the machine, is provided at an edge of the interchangeable sorting plate which edge extends in the longitudinal extension (long side). The holding strap in the holding recess may be able to be pressed against the carrier plate in a non-positive manner, e.g., by means of a retaining spring. The holding strap and the fixing element may be arranged opposite one another directly or with an offset relative to the transverse extension of the interchangeable sorting plate. The holding strap points in the direction of installation of the interchangeable sorting plate.\nIt is also possible to provide at least one stop element, which is a rigid part of the machine and with which the interchangeable sorting plate comes into contact in the direction of installation or at right angles to the longitudinal extension of the interchangeable sorting plate in the course of the insertion with the long side.\nIt is obvious that a plurality of adjusting surfaces, holding straps and/or fixing elements, with respective associated components, which are rigid parts of the machine, may also be provided within the framework of the present invention.\nThe present invention also pertains to a set of interchangeable sorting plates for a device according to the present invention, where the interchangeable sorting plates of a set of interchangeable sorting plates are designed for different sets of coins. This means that an interchangeable sorting plate is provided with sorting openings of the cross sections d depending on the diameters of the coins of the particular set of coins to be sorted. For example, a device used in Germany may have a set of interchangeable sorting plates comprising an interchangeable sorting plate for DM and an interchangeable sorting plate for euro. Using such a set of interchangeable sorting plates, the changeover to the other currency can be carried out rather easily and without any major installation and/or service effort especially also during the phase of transition. Therefore, the present invention also pertains to a process for converting a device according to the present invention to different sets of coins, in which case an interchangeable sorting plate for one set of coins is replaced with an interchangeable sorting plate for another set of coins.\nThe various features of novelty which characterize the invention are pointed out with particularity in the claims annexed to and forming a part of this disclosure. For a better understanding of the invention, its operating advantages and specific objects attained by its uses, reference is made to the accompanying drawings and descriptive matter in which preferred embodiments of the invention are illustrated."}
{"text": "German Patent Application No. 28 29 057, discusses a fuel-injection system for supplying fuel to a mixture-compressing internal combustion engine having external ignition as a function of operating parameters. The fuel-injection system includes a metal fuel-distributor line, which, via at least one branch line, is connected to at least one fuel injector, the branch line being embodied as a metal tube and connected to the fuel injector by way of a threaded connection. Easily bendable metal is used as material for the branch line. Situated between the threaded connection on the branch line and the fuel injector are thin-walled metal bellows, which compensate for a lateral offset between the beginning point of the branch line on the fuel-distributor line and the fitting position of the fuel injector, and which damp the operating noises originating in the fuel injector.\nDisadvantageous in the fuel-injection system from German Patent Application No. 28 29 057 is that the screw connection is not secured against an automatic loosening. Due to the vibrations of the internal combustion engine during operation, a screwed connection rigidly connected to the internal combustion engine is at a higher risk of coming unscrewed.\nHowever, especially in the case of directly injecting fuel injectors and at the high pressures required in this context, the screwed connection is safety-relevant and should not come unscrewed under any circumstances. The known fuel-injection system provides no suggestion for securing a screwed connection."}
{"text": "1. Field of the Invention\nThe present invention relates to a solid electrolytic capacitor and a method of manufacturing the same.\n2. Description of the Related Art\nIn general, solid electrolytic capacitors having a function of storing electricity are electronic components used for blocking direct currents and passing alternating currents. Among the solid electrolytic capacitors, the most representative one is a tantalum capacitor which is used in application circuits, of which the rate voltage range is low, as well as general industrial equipments. In particular, the tantalum capacitor is frequently used in circuits requiring an excellent frequency characteristic or for reducing noise of communication equipments.\nFIG. 1 is a perspective view of a conventional solid electrolytic capacitor, and FIG. 2 is a cross-sectional view of the conventional solid electrolytic capacitor. As shown in FIGS. 1 and 2, the solid electrolytic capacitor 10 includes a capacitor element 11 which determines the capacity and characteristic of the capacitor and is formed of dielectric ceramic powder, anode and cathode lead frames 13 and 14 which are connected to the capacitor element 11 so as to be easily mounted on a printed circuit board (PCB), and an epoxy case 15 which is molded of epoxy so as to protect the capacitor element 11 from the external environment and to form the shape of the capacitor element.\nIn one side of the capacitor element 11, a rod-shaped anode wire 12 is formed to project with a predetermined length.\nThe anode wire 12 has a planar surface 12a provided thereon, the pressed surface 12a increasing a contact area with the anode lead frame 13 and preventing the anode lead frame 12 from rocking from side to side during welding.\nThe capacitor element 11 is manufactured by the following process. First, dielectric ceramic powder is molded in a rectangular parallelepiped shape in a pressing process and is then sintered. Further, a dielectric oxide film is formed on the surface of the sintered body. Then, the body is dipped into a manganese nitrate solution such that a manganese dioxide layer composed of a solid electrolyte is formed on the outer surface of the body.\nA process of connecting the anode and cathode lead frames 13 and 14 to the capacitor element 11 manufactured in such a manner includes two steps. In the first step, the plate-shaped anode lead frame 13 is welded on the planar surface 12a of the rod-shaped anode wire 12, which projects from one side surface of the capacitor element 11 at a predetermined length, so as to derive an anode terminal. In the second step, a cathode terminal is derived through the surface of the capacitor element 11 or a conductive adhesive coated on the cathode lead frame 14.\nThen, the capacitor element 11 is electrically connected to the anode and cathode lead frames 13 and 14, respectively, and the epoxy case 15 is molded of epoxy. Then, the solid electrolytic capacitor is completed through a subsequent assembling process.\nThe above-described conventional electrolytic capacitor has the following problems.\nWhile the anode wire 12 and the anode lead frame 13 are directly welded, high-temperature heat is generated. The generated heat has an effect upon the capacitor element 11 through the anode wire 12, thereby damaging the capacitor element 11 which is vulnerable to heat.\nFurther, dielectrics are destroyed by the heat shock applied to the capacitor element 11 such that a product quality is degraded and defects occur. Therefore, a manufacturing cost increases.\nFurther, the anode lead frame 13 and the cathode lead frame 14 occupy such a large space in the epoxy case 15. Therefore, the capacitor element 11 is inevitably reduced in size within the epoxy case 15. As a result, the capacitance of the capacitor decreases."}
{"text": "Field of the Invention\nEmbodiments of the invention provide techniques for extracting pixel-level micro-features from image data. More specifically, embodiments of the invention relate to techniques for producing a micro-feature vector for the image data that is used to classify objects depicted in the image data.\nDescription of the Related Art\nSome currently available video surveillance systems provide simple object recognition capabilities. For example, a video surveillance system may be configured to classify a group of pixels (referred to as a “blob”) in a given frame as being a particular object (e.g., a person or vehicle). Once identified, a “blob” may be tracked from frame-to-frame in order to follow the “blob” moving through the scene over time, e.g., a person walking across the field of vision of a video surveillance camera. Further, such systems may be configured to determine the type of object that the “blob” is.\nHowever, such surveillance systems typically require that the objects which may be recognized by the system to be defined in advance. Thus, in practice, these systems rely on predefined definitions for objects to evaluate a video sequence. In other words, unless the underlying system includes a description for a particular object, i.e., has been trained, the system is generally incapable of recognizing that type of object. This results in surveillance systems with recognition capabilities that are labor intensive and prohibitively costly to maintain or adapt for different specialized applications. Accordingly, currently available video surveillance systems are often unable to identify objects, events, behaviors, or patterns as being “normal” or “abnormal' by observing what happens in the scene over time; instead, such systems rely on static object definitions.\nFurther, the static patterns recognized by available video surveillance systems are frequently either under inclusive (i.e., the pattern is too specific to recognize many instances of a given object) or over inclusive (i.e., the pattern is general enough to trigger many false positives). In some cases, the sensitivity of may be adjusted to help improve the recognition process, however, this approach fundamentally relies on the ability of the system to recognize predefined patterns for objects. As a result, by restricting the range of objects that a system may recognize using a predefined set of patterns, many available video surveillance systems have been of limited (on simply highly specialized) usefulness."}
{"text": "1. Field of the Invention\nThe present embodiments relate to a plasma display panel (PDP) used for a plasma display device, and more particularly, to a PDP that improves both contrast and discharge efficiency by reducing residual images and reflection of external light.\n2. Description of the Related Art\nFIG. 1 is a fragmentary perspective view of a conventional PDP as disclosed in Japanese Patent Laid-Open Publication No. 1998-172442.\nAs shown in FIG. 1, the conventional PDP includes an upper panel 1 and a lower panel 2, which are bonded to each other, and a space defined by the upper and lower panels 1 and 2 is filled with a discharge gas. The upper panel 1 includes a front substrate 60, pairs of sustain discharge electrodes 84, and a front dielectric layer 80. Each of the pairs of sustain discharge electrodes 84 include a Y-electrode 83 and an X-electrode 82, which are disposed on a bottom surface 60a of the front substrate 60. The pairs of sustain discharge electrodes 84 are covered with the front dielectric layer 80. The front dielectric layer 80 may be covered with a protective layer 90 that is typically formed of MgO. Meanwhile, the Y-electrode 83 includes a first transparent electrode 83b, which is formed of indium tin oxide (ITO), and a first bus electrode 83a, which prevents a voltage drop in the first transparent electrode 83b. Like the Y-electrode 83, the X-electrode 82 includes a second transparent electrode 82b and a second bus electrode 82a. \nThe lower panel 2 includes a rear substrate 10, address electrodes 20, a rear dielectric layer 30, barrier ribs 40, and phosphor layers 50. The address electrodes 20 are disposed on a top surface of the rear substrate 10 and intersect the pairs of sustain discharge electrodes 84. The address electrodes 20 are covered with the rear dielectric layer 30. The barrier ribs 40 are disposed on the rear dielectric layer 30 and define discharge cells along with the pairs of sustain discharge electrodes 84. The phosphor layers 50 are coated on inner surfaces of the discharge cells.\nFor the above-described conventional PDP, a discharge cell is selected due to address discharge induced between the address electrode 20 and the Y-electrode 83. Then, sustain discharge arises between the X-electrode 82 and the Y-electrode 83 of the selected discharge cell, so that the discharge cell emits light. More specifically, as the sustain discharge happens, the discharge gas filled in the discharge cell emits ultraviolet rays, which lead the phosphor layers 50 to emit red, green, and blue visible rays, thus producing the image of the PDP.\nIn the conventional PDP having the above-described construction, external light is projected on a front surface of a display region of the PDP. Owing to this reflection of external light, the PDP has higher reflection brightness and a lower contrast.\nTherefore, it has become increasingly necessary to cut off the reflection of external light in order to minimize an elevation in the reflection brightness of the PDP and improve the contrast thereof.\nFIG. 3 is a diagram for schematically explaining the principles of additive mixture of colors and subtractive mixture of colors.\nAs shown in FIG. 3, when the three primary colors of light, that is, red, green, and blue colors, are mixed using the additive mixture of colors, a central portion in which all the three colors are mixed emits white light, and portions in which two of the three colors are mixed emit yellow, azure, and violet light, respectively. Here, it is said that two colors disposed opposite each other across a white color are complementary to each other. In mixing the complementary colors, the mixed color lowers in both brightness and chroma and, preferably, becomes a black color.\nWhen yellow, azure, and violet colors are mixed using the subtractive mixture of colors, a central portion in which all the three colors are mixed takes on a black color, and portions in which two of the three colors are mixed take on red, green, and blue colors, respectively. For example, when the red and azure colors, which are complementary to each other, are mixed, the mixed color lowers in both brightness and chroma. Therefore, by making use of complementary colors as described above, it is expected that there will be improvements in contrast and other problems."}
{"text": "The present invention relates to Ethernet switches and switching, and more particularly, this invention relates to providing I2C bus over Ethernet.\nA significant portion of network switches have various devices which utilize I2C bus as an interface to a local host external of the network switch or a local processor of the network switch, such as a central processing unit (CPU). Devices which communicate using an I2C bus will be referred to throughout this document as I2C bus devices. These I2C bus devices are typically controlled by the local processor over the I2C bus of the network switch. Some typical I2C bus devices include, but are not limited to, fan controllers, power supply units, various thermal gauges, serial electrically erasable programmable read only memory (SEEPROM) memories of various subunits, etc. In cell-based distributed fabric protocol (cDFP) and/or Ethernet-based distributed fabric protocol (eDFP) scaled out systems (referred to herein collectively as DFP systems), there are multiple individual system members which, in order to function properly, need to be controlled by the system master in real time. However, in conventional distributed systems or clusters, there is no way to provide real-time I2C bus over Ethernet."}
{"text": "1. Technical Field\nThis invention relates to internal combustion engines and in particular engines using a magnetorheological lubricant.\n2. Background Art\nIt is generally desirable to reduce the fuel consumption of an engine to reduce the running costs of the engine and the emissions from the engine. The emissions from an engine are closely related to the volume of fuel consumed by the engine and this is particularly so in the case of CO2 emissions. Two major factors affecting fuel economy are: 1) Cold Start Cranking (CSC); and 2) Engine Warm Up (EWU).\nFor CSC, the starter motor is normally required to rotate an engine with cold, thick, viscous engine oil combined with un-lubricated, and hence, high-friction internal engine components. Both the aforementioned issues provide an increased inertia for the starter motor directly putting an additional drain upon the vehicle's battery. The battery charge is replenished via the alternator, when the engine has started, but at a cost to fuel economy.\nFor EWU, once the engine has started, the engine is once again required to work harder to overcome the cold and viscous engine oil until optimum engine operating temperatures and related low oil viscosity is achieved. During EWU, extra fuel is consumed to compensate for higher frictional losses at a further cost to fuel economy.\nIn a traditional engine, the effects of thick, viscous oil and un-lubricated internal engine components are not uniform throughout the engine. For example, the crankshaft main bearings will have a higher and significant friction drain when compared to the camshaft bearings.\nReducing the viscosity of a conventional lubricant used to lubricate an engine reduces fuel consumption due to reduced frictional losses and a reduction in the power required to pump the lubricant through the engine. However, lubricants having a very low viscosity at ambient temperature have even lower viscosity as temperature increases. Although this reduction in viscosity with temperature is beneficial to fuel economy, there is a limit to how low the viscosity can go before the lubricating film breaks down at high temperatures and load with potentially serious consequences such as increased bearing wear and very high levels of friction. It is therefore normal practice to use a lubricant that is sufficiently viscous at high temperature to resist film breakdown and accept the increased and undesirable higher friction and high pumping losses at lower temperatures that will result from the use of oil that is more viscous than is desirable at ambient temperature."}
{"text": "Use of imaging sensors in vehicle imaging systems is common and known. Examples of such known systems are described in U.S. Pat. Nos. 5,949,331; 5,670,935 and/or 5,550,677, which are hereby incorporated herein by reference in their entireties."}
{"text": "This invention relates to a novel class of polypropylenes that possess elastic character and are hybrid polymers having steroregularity intermediate between those that are highly syndiotactic (or highly isotactic) and those that are atactic, as well as catalysts that produce these hybrid polymers.\nElastic polypropylene was first isolated by solvent extraction of polypropylene obtained by Ziegler-Natta polymerization of propylene (inter alia, U.S. Pat. No. 3,175,999). A process for preparing primarily isotactic, fractionable, elastomeric polyolefins has been described (U.S. Pat. No. 4,335,225), using a catalyst prepared by reacting a solid support such as alumina with, e.g., a tetraalkylzirconium compound. U.S. Pat. No. 5,594,080 discloses preparation of a stereoblock elastomeric polypropylene using a metallocene catalyst containing two unbridged phenylindenyl ligands connected to zirconium, i.e. (phenylindenyl)2ZrCl2. Chien et al. (J. Am. Chem. Soc., 112, 2030 (1990); Macromolecules, 25, 1242 (1992); Macromolecules, 24, 850 (1991)) prepared a thermoplastic, elastomeric polypropylene containing about 30% crystalline isotactic polypropylene, using a metallocene catalyst.\nCollins et al. (Macromolecules, 28, 3771 (1995); Macromol. Symp., 98, 223 (1995)) describe the use of zirconium and hafnium metallocene catalysts to prepare elastic polypropylene. These polymers had 29-54% mmmm pentads, i.e. they were predominantly isotactic. Significantly, elastic polypropylene could only be obtained when the molecular weight was greater than 50,000 and when the mmmm content of the polymer exceeded 38%, using exclusively hafnium complexes.\nEuropean Patent Application No. 666,267 describes the use of a metallocene catalyst of the type {fluorenyl-C2H4-indenyl}ZrCl2 to prepare polypropylene in the presence of hydrogen. Polymer properties are not described. The same catalyst was used (Organometallics, 13, 647 (1994)) to polymerize propylene in toluene in the absence of hydrogen; highly isotactic polymers having mmmm contents ranging from 38 to 64% and melting points of 104-110xc2x0 C. were obtained. When the C2H4 bridging group was replaced by xe2x80x94SiMe2xe2x80x94 (U.S. Pat. No. 5,391,789), only a viscous, tacky oil was obtained.\nEuropean Patent Application No. EP 707,016 describes a catalyst system comprising {2-Me-4-naphthyl-indenyl-SiMe2-fluorenyl}ZrCl2, a Lewis acid compound and an organoaluminum compound, to prepare elastomeric polypropylene. U.S. Pat. No. 5,516,848 describes a process for making elastomeric polypropylenes employing a mixture of two different metallocene catalysts, one of which was a monocyclopentadienyl transition metal compound.\nMetallocene catalysts generally can be defined as catalysts comprising one or more cyclopentadienyl groups, including substituted cyclopentadienyl groups, in combination with a transition metal, typically a metal of Group IVB of the Periodic Chart (i.e., titanium, zirconium, and hafnium). Useful metallocenes have been described comprising two cyclopentadienyl-type rings, in which both rings are identical (symmetrical metallocenes) and in which the two rings are different, due either to differing substitution patterns on identical ring systems or to the presence of two different ring systems (unsymmetrical metallocenes). The first reports of the use of cyclopentadienyl metallocenes to catalyze the polymerization of 1-olefins appeared in 1957. British Patent Application GB 934,281 first described metallocene catalysts of the type (R-indenyl)2MX2 and (R-fluorenyl)2MX2, where R represents an alkyl or aryl substitutent on the aromatic ring, M represents titanium, zirconium or hafnium, and X represents a halide or an alkyl or alkoxy group, or the like, having from 1 to 12 carbon atoms. Since then, numerous patents and other publications have provided ample evidence of the commercial interest in these materials.\nFour limiting cases illustrate the possibilities for stereoregularity in polypropylene. The isotactic structure is typically described as having the methyl side groups oriented so that they are all on the same side of the plane containing the polymer backbone. Another way of describing the isotactic structure employs Bovey\"\"s NMR (nuclear magnetic resonance) nomenclature: An isotactic diad, represented by m (for meso), is a pair of propylene units placed in the polymer chain so that both methyl side groups lie on the same side of the plane containing the polymer backbone. Similarly, a mmmm pentad represents five monomer units with their methyl groups all oriented the same way. 13C NMR can be used to determine the relative amounts of the various pentads present in a propylene polymer.\nIn contrast, in the syndiotactic structure, the methyl groups are oriented alternately above and below the plane containing the polymer backbone. A pair of propylene units with this up-down arrangement comprises an r (for racemic) diad, and rrrr would correspond to five similarly arranged monomer units.\nHemi-isotactic polypropylene is similar to the isotactic kind except that every other methyl side group has a random orientation.\nFinally, in the atactic polymer, the orientation of the methyl groups with respect to the plane containing the polymer backbone is random and the number of r and m diads is the same.\nThe physical properties of polypropylene depend greatly on polymer stereoregularity. For example, isotactic and syndiotactic polypropylene are crystalline polymers that are insoluble in hydrocarbon solvents such as cold xylene. Because of the presence of crystallites that scatter light, they are often opaque or translucent. Additionally, these polymers are stiff, inelastic, and high-melting, their melting points being 171 and 138xc2x0 C., respectively.\nIn contrast, hemi-isotactic and atactic polymers are non-crystalline and soluble in cold xylene. Atactic polypropylene is transparent, flexible and elastomeric.\nBecause isotactic and syndiotactic polypropylenes are such useful materials, their preparation has been the object of extensive research. For example, catalysts that produce isotactic polypropylene are disclosed in U.S. Pat. Nos. 4,794,096 and 4,975,403; catalysts that lead to syndiotactic polyolefins are described in U.S. Pat. Nos. 3,258,455, 3,305,538, 3,364,190, 4,892,851, 5,155,080 and 5,225,500.\nThe precise effect of changes in catalyst structure on catalyst activity and polymer stereoregularity is difficult to predict. For example, U.S. Pat. No. 5,459,218 describes catalysts of the type {flu-bridge-Cp}ZrCl2, wherein flu=fluorenyl, bridge=Me2Si or Ph2Si, and wherein Me=methyl, Ph=phenyl, and Cp=cyclopentadienyl. It is broadly stated and claimed that the propylene polymers contain at least 50% rr content, whereas the syndiotactic content actually achieved ranged from 74 to 82%. The highest molecular weight given for the propylene polymers is 66,000. No unusual properties of these polymers were disclosed nor were any clues provided as to how the catalysts could be modified to produce polymers having a lesser or greater syndiotactic content. By comparison, U.S. Pat. No. 4,892,851 describes a catalyst of a very similar type, namely {flu-CMe2-Cp}ZrCl2, which is reported to produce even more highly syndiotactic (92% rr) polypropylene.\nA relationship between catalyst symmetry and polypropylene stereostructure has been proposed (Trends in Polymer Science, 2, 158 (1994)). Symmetry elements, typically described in terms of xe2x80x9cpoint groupsxe2x80x9d of a bis(fluorenyl)MX2-type catalyst are shown in FIGS. 1 and 2 (comparative). FIG. 1 illustrates the different symmetry elements that can occur in metallocene catalysts and FIG. 2 provides illustrative examples. A catalyst (FIG. 1 and, e.g., FIG. 2 entry 20) having point group C2v has three symmetry elements: the structure is unchanged after a 180 degree rotation about the bisector of the MX2 plane and there are two mirror planes of symmetry, reflections in which leave the structure unchanged. One contains the MX2 plane \"sgr\"h and the other is orthogonal to and bisects the MX2 plane \"sgr\"v. Looking straight on at the MX2 plane, one sees that the top and bottom as well as the left and right hand side of the catalyst molecule are the same. A catalyst 21 having point group C2 (entry 21 in FIG. 2) contains only a C2 axis: the top and bottom of the molecule are the same, so that 180 degree rotation about the bisector of the MX2 plane leaves the structure unchanged. Catalyst 22 having point group CS contains only a mirror plane of symmetry orthogonal to and bisecting the MX2 plane. Viewed as described above, the left- and right hand sides of the catalyst are the same but the top and bottom are different. Catalyst 25 of C1 point group symmetry contain none of these three symmetry elements.\nAccording to FIG. 2, catalysts of point group C2v are said to produce atactic polypropylene, except that Cp2TiPh2 (24) atypically produces a stereoblock isotactic polymer (perhaps by a different polymerization mechanism) at low temperatures. Metallocenes having C2 and CS point group symmetry are said to produce isotactic and syndiotactic polymers, respectively. However, U.S. Pat. No. 4,769,510 describes catalysts of the type {ind-bridge-ind}MX2, which exists in two isomeric forms. One form is chiral and consists of a d,l pair of enantiomers; it produces isotactic polypropylene. The second (meso) form is achiral and produces atactic polypropylene (Organometallics, 14, 1256 (1995)). Catalysts with C1 symmetry are said to produce either hemi-isotactic 23 (in FIG. 2) or stereoblock isotactic-atactic 25 (in FIG. 2) polypropylene.\nEuropean Patent Application No. 537,130 describes the effect of the size of an organic group attached to a cyclopentadienyl ring in metallocenes of the type {Cp-CMe2-flu}ZrCl2 (Cs point group symmetry), which was shown to produce syndiotactic polypropylene. Introduction of a methyl group into the 3-position of the cyclopentadienyl ring led to {3-MeCp-CMe2-flu}ZrCl2 (C1 point group symmetry), which produced a hemi-isotactic polymer. When the 3-CH3 group was replaced by a t-butyl group to form {3-t-BuCp-CMe2-flu}ZrCl2 (also C1 point group), isotactic polypropylene was obtained instead.\nU.S. Pat. No. 5,459,218 relates to syndiotactic polypropylene prepared using silyl bridged metallocenes and having molecular weights up to 66,000. U.S. Pat. No. 5,668,230 discloses certain specific ethylene bridged fluorenyl-containing metallocenes to catalyze olefins polymerization.\nOne skilled in the art can only conclude that there are no reliable, general correlations between the structure of a metallocene catalyst and the stereospecificity or the stereoregularity of polymers produced thereby.\nBriefly, the present invention provides propylene homopolymers that are elastomeric and are soluble in at least one nonpolar organic solvent selected from the group consisting of toluene, xylene, heptane, and hexane, the polymer comprising\na) greater than 3 weight percent and up to 45 weight percent of homotactic sequences, each having only r or m diads, all of which homotactic sequences have a helical length in the range of about 20 to 150 xc3x85, and\nb) in the range of 55 to 97 weight percent of the total composition being the sum of\n1) homotactic sequences, each having only r or m diads, and being less than 20 xc3x85 in helical length and having fewer than 10 repeat units with mmmm pentads being present in the range of 0 to 35 (preferably 0 to 31) weight percent of the total composition, and\n2) heterotactic sequences having unequal numbers of r and m diads,\nthe polymer having a weight average molecular weight (Mw) of at least about 70,000, preferably greater than 70,000 and up to 2,000,000, more preferably 75,000 to 1,000,000, and most preferably 80,000 to 500,000. The two b) components can be present in any amount so long as their sum amounts to 55 to 97 weight percent of the total composition.\nThe propylenes of the invention include homotactic sequences of length as specified above having all r or all m diads. As is appreciated by those skilled in the art, these homotactic sequences can assume a helical configuration. Connecting these homotactic sequences are similar homotactic sequences of less than 20 xc3x85 in helical length and having less than 10 repeat units as well as heterotactic sequences having both r and m diads present in unequal numbers. These polymers exhibit elastomeric character that varies from a stiff rubber to that of a very stretchy rubber band. The homotactic sequences can be all r or all m and preferably they are all r, i.e., syndiotactic. In general, the polymers comprise discontinuous small (20 xc3x85 to 150 xc3x85 in length) hard segments in a continuous softer matrix.\nThe length of the homotactic sequences can be estimated in the following way. Helical chains of polypropylene run approximately parallel to the c-axis of the unit cell of orthorhombic syndiotactic polypropylene. Four monomer units are contained in the unit cell which is 7.6 xc3x85 long. Therefore, each molecule of propylene contributes 7.6/4 or 1.9 xc3x85 to the chain length. This is a reasonable estimate because if the chain were fully extended, each propylene would contribute the distance of two Cxe2x80x94C bonds or 3 xc3x85 to the chain length.\nThe heterotactic sequences comprise both r and m diads and are of variable lengths and are randomly dispersed as to size and distribution in the polymer chain.\nPreferably, the propylene polymers of the invention also have one or more of the following characteristics:\ni) a stereoregularity index between 1.30 and 10.0, preferably 1.30 to 7.00 and more preferably 1.60 to 6.40.\nii) a heat of fusion (xcex94Hfus) that is less than 50% of the xcex94Hfus of 100% isotactic polypropylene when mm greater than rr; or less than 50% of the xcex94Hfus of 100% syndiotactic polypropylene when rr greater than mm, (xcex94Hfus values are as given in J. Brandrup et al., POLYMER HANDBOOK, 3d Edition, John Wiley and Sons, NY (1989), section V, p. 29); and\niii) optical clarity.\nIn another aspect, this invention provides blends of two or more different propylene polymers, at least one of which exhibits the properties described above. Blends of one or more propylene polymers as described above can further comprise crystalline or amorphous polypropylene, tackifying resins, antioxidants, fillers, and other adjuvants known in the art.\nIn yet another aspect, this invention provides copolymers of propylene wherein the propylene sequences are as defined above, and wherein the comonomers can be olefins having 2 to 20, preferably 2 to 8, carbon atoms. Copolymers can be blended with one another as well as with homopolymers which provides a way to tailor or modify polypropylene properties. Molecular weight of copolymers (Mw) can be in the range of 35,000 to 2,000,000, preferably 50,000 to 1,000,000, more preferably 70,000 to 500,000.\nIn a further aspect, this invention provides metallocene catalysts for propylene polymerization, having the structure\n{ligand1-bridge-ligand2}MX2, \nwherein\nligand1 and ligand2 are different and are selected from the group consisting of substituted and unsubstituted cyclopentadienyl (Cp), indenyl (ind), fluorenyl (flu), 4,5-dihydrocyclopentaphenanthryl (H2CPA), and cyclopentaphenanthryl (CPA) ring groups, wherein, when present, ring group substituents can be selected from the group consisting of\nii) C6-C20 aryl,\niii) C7-C20 alkylaryl,\niv) C4-C7 cycloalkyl,\nv) (xe2x80x94CH2xe2x80x94)n, wherein n is 2, 3, 4, or 5, or (xe2x80x94CHxe2x95x90CHxe2x80x94)m wherein m is 1, 2, 3, or 4, connecting two positions (i.e., adjacent or non-adjacent ring carbon atoms) in the same ring structure, preferably the 4 and 5 positions in fluorenyl or indenyl,\nvi) fused aromatic rings, and\nvii) fused aromatic rings substituted by any one of groups i)-v);\nbridge is a linking group joining ligand1 and ligand2 at C-1 of Cp or ind ligands or C-9 of flu and CPA ligands and is selected from the group consisting of\ni)  greater than CR1R2,\nii)  greater than SiR1R2,\niii) xe2x80x94CR1R2xe2x80x94CR3R4xe2x80x94,\niv) xe2x80x94SiR1R2xe2x80x94SiR3R4xe2x80x94,\nv) xe2x80x94CR1R2xe2x80x94SiR3R4xe2x80x94,\nwherein R1, R2, R3, and R4 can be the same or different and are selected from the group consisting of H, C1-C20 straight-chain or branched alkyl groups, C6-C20 aryl groups, and C3-C8 cycloalkyl groups;\nM is a metal atom selected from the group consisting of Zr, Hf, and Ti, preferably is zirconium or hafnium, and most preferably zirconium, and\nX is selected from the group consisting of Cl, Br, I, C1-C20 straight-chain or branched alkyl groups, C6-C20 aryl groups, C7-C20 alkaryl groups, and C7-C20 aralkyl groups,\nwherein the metallocene catalyst exhibits an asymmetry parameter of 1.03 to 1.69, preferably 1.03 to 1.45, more preferably 1.05 to 1.35,\nwith the proviso that the metallocene catalyst provides propylene polymers described above.\nIn a still further aspect, this invention also provides a method of making novel catalyst precursor compounds of the type\nH{ligand1-CH2CH2-ligand2}H \nwherein ligand1 and ligand2 are as previously defined,\nthe method comprising the steps:\n1) deprotonating the hydroxyl group of a compound of the type\nligand1-CH2CH2OH; \n2) reacting the deprotonated compound with a perfluoroalkylsulfonyl fluoride of the type RfSO2F to obtain a stable perfluoroalkylsulfonate of the type\nligand1-CH2CH2xe2x80x94OSO2Rf; \nwherein Rf means an alkyl group having 1-20 carbon atoms or an aryl group having 5 to 12 carbon atoms in which at least 75% of H atoms have been replaced by F atoms;\n3) condensing the perfluoroalkylsulfonate with the conjugate base of a compound of the type\nligand2-H; \nwherein ligand1 and ligand2 are as previously defined and Rf is independently selected from the group consisting of highly fluorinated or perfluorinated alkyl radicals containing from 1-20 carbon atoms.\nIn yet a further aspect, there is provided a novel class of metallocene catalysts that are particularly useful to prepare hybrid polymers of the present invention, i.e. those having intermediate stereoregularity between those that are highly syndiotatic (or highly isotactic) and those that are atactic, the novel catalysts of the present invention having the formula: \nwherein\nstrap can be a (xe2x80x94CH2xe2x80x94)n group or a (xe2x80x94CHxe2x95x90CHxe2x80x94)m group,\nbridge is as previously defined,\nMa is Zr or Hf,\nn is 1, 2, 3, or 4 and m is 1, 2, 3, or 4, and\nX is as previously defined,\nwith the proviso that when Ma=Zr, X=Cl, and bridge=C2H4, then strap can (xe2x80x94CH2xe2x80x94)n or (xe2x80x94CHxe2x95x90CHxe2x80x94)m, wherein n=3, 4 or 5, and m=2, 3, or 4, preferably m is 2, 3, or 4 and can be designated mxe2x80x2.\nIn yet another aspect, this invention also provides a method of controlling the stereostructure of polypropylene by choosing the shape of a metallocene catalyst, as defined by its asymmetry parameter, from the class of catalysts having the structure\n{ligand1-bridge-ligand2}MX2, \nwherein\nligand1 and ligand2, bridge, M, and X are as previously defined, and wherein the metallocene catalyst exhibits an asymmetry parameter of 1.03 to 1.69, and the stereoregularity index of said polypropylene increases from about 1.30 to about 10.00 with an increase in the asymmetry parameter of the catalyst.\nPreparation of the hybrid polypropylenes of the invention is accomplished by combining propylene with a metallocene catalyst, preferably in the absence of solvent for the monomer, in an inert atmosphere at a pressure in the range of 69 to 6890 Kpa at a temperature in the range of xe2x88x9220xc2x0 C. to about 120xc2x0 C. When a solvent is used, it can be a hydrocarbon, preferably aliphatic, cycloaliphatic, or aromatic, more preferably it is toluene or cyclohexane. Activation of the metallocene is accomplished by activators known in the art, preferably methylaluminoxane alone or in combination with trialkylaluminum compounds such as trimethylaluminum, triethylaluminum, or triisobutylaluminum.\nAs used herein:\nxe2x80x9casymmetry parameter (A.P.)xe2x80x9d is defined as the ratio of the van der Waals surface area of the larger ligand to that of the smaller ligand, wherein the van der Waals surface area can be calculated by CAChe(trademark) Satellite molecular modeling program (version 3.8) (Oxford Molecular Ltd., Oxford, United Kingdom) on a Macintosh computer; this is a measure of the asymmetric character of a molecule;\nxe2x80x9catacticxe2x80x9d means xe2x80x9cpolypropylene having no dominant tacticity and having 25% mm triads and 25% rr triads, as determined by NMR spectroscopy;\nxe2x80x9cCp or cpxe2x80x9d means cyclopentadienyl;\nxe2x80x9cCPAxe2x80x9d means cyclopentaphenanthryl;\nxe2x80x9celastomericxe2x80x9d or xe2x80x9celasticxe2x80x9d means a material that at room temperature can be stretched repeatedly to at least twice its original length and, immediately upon release of the stress, returns with force to its approximate original length;\nxe2x80x9centanglement molecular weight ({overscore (M)}e)xe2x80x9d means average molecular weight of chains between entanglements in a polymer;\nxe2x80x9cfilmxe2x80x9d means a self-supporting layer;\nxe2x80x9cfluxe2x80x9d means fluorenyl;\nxe2x80x9cgroupxe2x80x9d or xe2x80x9cradicalxe2x80x9d or xe2x80x9ccompoundxe2x80x9d or xe2x80x9cligandxe2x80x9d or xe2x80x9cmonomerxe2x80x9d or xe2x80x9cpolymerxe2x80x9d means a chemical species that allows for substitution or which may be substituted by conventional substituents which do not interfere with the desired product; e.g., substituents can be alkyl, aryl, phenyl, etc.; and\nxe2x80x9cH2CPAxe2x80x9d means 4,5-dihydrocyclopentaphenanthryl;\nxe2x80x9chelical lengthxe2x80x9d means the length of a segment in a coiled configuration in contrast to the filly extended polymer segment;\nxe2x80x9cheterotacticxe2x80x9d sequences are those having both r or m diads;\nxe2x80x9chighly fluorinatedxe2x80x9d (Rf) means having at least 75% of H atoms in an alkyl cycloalkyl, or aryl group substituted by fluorine atoms;\nxe2x80x9chighly isotacticxe2x80x9d means polypropylene having an mm content of at least 80%, as determined by NMR spectroscopy;\nxe2x80x9chighly syndiotacticxe2x80x9d means polypropylene having an rr content of at least 80%, as determined by NMR spectroscopy;\nxe2x80x9chomotacticxe2x80x9d sequences are those having diads which are all r or all m.\nxe2x80x9cindxe2x80x9d means indenyl;\nxe2x80x9cisotacticxe2x80x9d or xe2x80x9cisotactic-richxe2x80x9d means polypropylene having tacticity shown in FIG. 2 (second entry) and having an mm content of greater than the rr triad content, as determined by NMR spectroscopy;\nxe2x80x9cMexe2x80x9d means methyl;\nxe2x80x9cmetallocenexe2x80x9d means a metal-organic compound characterized by xcfx80-bonds between a transition metal and a cyclopentadienyl, indenyl, or fluorenyl ligand, or such substituted ligands;\nxe2x80x9cmm triadsxe2x80x9d refers to an isotactic structure having three propylene units placed in a polypropylene backbone such that all three methyl side groups lie on the same side of the plane containing the polymer backbone, as determined by NMR spectroscopy;\nxe2x80x9cmr triadsxe2x80x9d refers to a stereoregular polypropylene structure having three successive propylene units placed in a polypropylene backbone such that, relative to the first propylene unit, the methyl group of the second propylene lies on the same side of the plane containing the polymer backbone and the methyl group of the third propylene lies on the opposite side of the plane containing the polymer backbone, as determined by NMR spectroscopy;\nxe2x80x9cnanocystallinexe2x80x9d means having crystallites in the nanometer size range (largest length), preferably 2 to 200 nm;\nxe2x80x9coptical clarityxe2x80x9d means transmittance of greater than 80% of light of wavelengths of 400-750 nm;\nxe2x80x9cpeak molecular weight,xe2x80x9d Mp, means the maximum molecular weight in a curve relating molecular weight and the abundance of a species of a given molecular weight determined by gel phase (or size exclusion) chromatography;\nxe2x80x9cPhxe2x80x9d means phenyl;\nxe2x80x9crr triadsxe2x80x9d refers to a syndiotactic structure having three propylene units placed in a polypropylene backbone such that each successive methyl side group is alternately above and below the plane containing the polymer backbone, as determined by NMR spectroscopy;\nxe2x80x9csolublexe2x80x9d means dissolves to an extent of greater than 98 weight percent at a temperature up to the boiling point of the stated solvent;\nxe2x80x9cstereoregularity index (S.I.)xe2x80x9d of a propylene polymer is defined as the ratio of the percentage of mm triads to rr triads, wherein the ratio represents the larger of mm or rr over the smaller of mm or rr (i e., the ratio is positive and greater than 1); and\nxe2x80x9csyndiotacticxe2x80x9d or xe2x80x9csyndiotactic-richxe2x80x9d means polypropylene having tacticity shown in FIG. 2 (third entry) and having an rr content of greater than the mm triad content, as determined by NMR spectroscopy.\nMetallocene catalysts that produce highly isotactic or highly syndiotactic polymers exhibit a number of chemical and/or structural similarities:\n1) These catalysts very often contain a MX2 unit (M=Ti, Zr, or Hf, X=Cl, Br or CH3, most commonly Cl) bonded to two variously substituted cyclopentadienyl-type ligands such as cyclopentadienyl itself, 1-indenyl or 9-fluorenyl.\n2) The two cyclopentadienyl-type ligands may be connected by a bridging group joined at the C-1 position in the indenyl moiety or the C-9 position in the fluorenyl moiety. The groups xe2x80x94CH2xe2x80x94CH2xe2x80x94 and xe2x80x94SiMe2xe2x80x94, wherein Me=methyl, are common examples.\n3) During the course of activation by an organoaluminum compound such as methylaluminoxane or mixtures of methylaluminoxane and other co-catalysts, Cl is replaced by one or two alkyl groups derived from the organoaluminum compound so that equivalent results for compounds containing a MCl2 or M(Me2) moiety are often seen.\n4) Once a basic catalyst structural type that provides isotactic or syndiotactic polymer has been recognized, stereoregularity can be increased by further adjustments of the structure through introduction of organic groups on the cyclopentadienyl-type ligands.\nHowever, prior to the present invention it was not possible to predict the physical properties of polypropylene of intermediate stereoregularity based on the shape of the metallocene catalyst.\nMetallocene catalysts of the invention produce novel propylene polymers, herein sometimes referred to as hybrid polymers, whose stereoregularity lies between highly syndiotactic (or highly isotactic) and atactic. Polymers thus obtained exhibit novel and useful properties, for example resistance to creep at high temperatures or extreme extensibility, which properties have not been previously observed in polypropylene."}
{"text": "There is known a technology for detecting the occurrence of a misfire in each of cylinders of an engine. See, Japanese Laid-Open Patent Publication No. 2009-280082 (hereinafter referred to as “JP2009-280082A”). According to JP2009-280082A, when a management ECU (117) judges that a misfire has occurred in an internal combustion engine (107), the management ECU turns on a warning lamp (125) (see paragraph [0027]).\nThere also is known a technology for detecting an abnormal compression pressure in each of cylinders of an engine. See, Japanese Laid-Open Patent Publication No. 2004-019465 (hereinafter referred to as “JP2004-019465A”). According to JP2004-019465A, while a fuel system and an ignition system are inactivated, the engine is cranked in order to rotate a crankshaft (1a), and a rotational variation of the engine is detected using a difference between instantaneous rotational speeds at preset crankshaft angles in compression strokes of the cylinders. An abnormal compression pressure is detected based on the rotational variation (see claims 1 through 6)."}
{"text": "US 2007/0040533 discloses in its title an input waveform control in a switching power supply and discloses in its abstract a recognition that a filter size can be reduced substantially as a power factor is permitted to deviate below unity in systematic ways. Specific, computable waveforms permit the use of a minimum filter size, given a desired target power factor. US 2007/0040533 further discloses in its FIG. 8 an output voltage resulting from an input voltage and a predefined input current and further discloses in its paragraph 0043 that, for a converter having a 200 μF output capacitor, this output voltage shows a relatively small 120 Hz ripple. The output capacitor is responsible for reducing this ripple. So, in case the output capacitor has a decreased value, the ripple will get an increased value.\nThis prior art disclosure is disadvantageous owing to the fact that the ripple in the output voltage is still too large. When using the converter for supplying a light source, this ripple will result in visible flicker. The prior art disclosure is further disadvantageous owing to the fact that the converter uses an electrolytic output capacitor having a relatively large value. Such an electrolytic output capacitor has a relatively short life time, especially at higher temperatures."}
{"text": "The invention relates to rubber mixtures which crosslink with moisture and which comprise a niobium compound or tantalum compound and polymers having alkoxysilyl terminal groups.\nSilicone rubber mixtures which can be stored with exclusion of water and on ingress of water crosslink at room temperature to give elastomers have been known for a long time, and this particularly applies to single-component silicone rubber mixtures (RTV-1). Large amounts of these products are used by way of example as jointing compositions in the construction industry. RTV-1 systems generally comprise a polymeric, mostly linear siloxane, a compound having crosslinking action and comprising groups that are readily hydrolysable, a plasticizer, and optionally other additions, such as a crosslinking catalyst, processing aids, pigments, and fillers. Crosslinking catalysts mainly used currently are Sn catalysts, e.g. dibutyltin diacetate or dibutyltin laurate. However, the recent toxicological classification of these tin derivatives has resulted in increasing criticism of these compounds, and the amounts permitted for use have been reduced accordingly. The use of Sn catalysts is therefore forbidden in elastomers which come into contact with foods, pharmaceuticals, or prostheses. There is a possibility that Sn-containing catalyst systems may be banned completely in the medium or long term. There is therefore a requirement for systems which can operate without tin catalysts but which at the same time have hardening performance comparable with that of conventional tin systems.\nVarious Sn-free catalyst systems have already been disclosed in the patent literature, especially using Ti compounds and Zr compounds. By way of example, U.S. Pat. No. 3,334,067 describes the use of catalysts such as Ti(OiPr)2(acac)2 or diisopropoxytitanium bis(ethyl-acetoacetate) which have improved shelf life when compared with standard Sn catalysts. Other titanium chelate compounds are mentioned inter alia in U.S. Pat. No. 3,542,901. Disadvantages of these types of catalyst derive from yellowing and incompatibility with aminic coupling agents. EP 102,268 claims the use of titanium-siloxy compounds, such as Ti(OSiMe3)4 or Ti(OSiMe3)2(On-Bu)4 and, respectively, Zr alkoxides, such as Zr(On-Bu)4. Siloxy-substituted Zr compounds are also used as condensation catalysts in JP 2004-238617A. JP 2007-131799 A describes niobium- and tantalum-alkoxy compounds of the general formulaM(OR)5 M=Nb, Ta R=alkylwhich are extremely susceptible to hydrolysis.\nA feature common to all of said Sn-free compounds is that, when comparison is made with the familiar Sn systems, there are disadvantageous properties, for example coloring of the rubber compositions, or excessive susceptibility to hydrolysis, or inadequate shelf life.\nThe present invention is based on the object of providing Sn-free alkoxy-crosslinking rubber mixtures which are not toxicologically hazardous and which give sufficiently high speed, and tack-free hardening."}
{"text": "This invention relates generally to investment casting, and more specifically relates to water disruptible cores useful in investment casting procedures, and to the method of preparation of such cores.\nIn the well-known procedures applicable to investment casting, a precise wax pattern is initially prepared corresponding to the metal article intended for ultimate manufacture. The wax pattern is thereupon repeatedly coated or \"invested\" with particulate ceramic material, so as to build up a ceramic shell mold of a desired thickness. Since the article to be manufactured often includes complex openings, cavities and the like, one or more ceramic cores are commonly disposed with respect to the wax pattern, so that when the wax is later removed from the shell mold, the cores remain in place, to form the appropriate cavities and the like, when the final metal casting is prepared.\nSeveral techniques are known for removing the wax pattern from the thus formed shell mold. Among the most common and acceptable is the use of high pressure and high temperature steam, which is accomplished by subjecting the shell-encased wax pattern to autoclaving. Once the wax is thus removed, the shell mold is fired in a furnace to remove wax residue and moisture and then the shell mold carrying the core or cores is ready for use for casting of the metal article.\nUpon the metal casting being completed, the ceramic core or cores must be removed from the cast article. In many instances this can be readily assured by forming the core from a material which is soluble in caustic alkali. However, this procedure cannot be used where one desires to cast metals or alloys such as aluminum, because in these instances the metal or alloys would themselves be chemically attacked by the caustic alkali. Accordingly, where casting of aluminum and similar metals and alloys are involved, it has been known in the prior art to form the ceramic cores of a material which is soluble, or at least includes sufficient soluble portions as to be thoroughly disruptible by water, so that after the casting is prepared, the core or cores remaining with the casting may be removed by use of water. Prior commercially available cores, of this type however, have not been capable of satisfactorily withstanding the autoclaving conditions which are discussed above, and which constitute an intermediate step to which the core is subjected prior to its use in casting.\nPrior art disclosures of water soluble or disruptible cores of the general type discussed above, include, for example, U.S. Pat. No. 4,629,708. The core described in this patent is a refractory body based upon a water soluble salt and a calcium silicate, the latter being present in amount of at least 45% by weight. The water soluble salt and the calcium silicate have a specified particle size distribution, and the calcium silicate is preferably a wollastonite. The patent discloses that this core may be removed from the cast or molded article by dissolution in water. The said patent also states that one can enhance the resistance of bodies made according to same to adventitious water, e.g., during storage or transport, by treating the cores after forming with a waterproofing material material. The cores may thus be coated and slightly impregnated by dipping them in an unsaturated polyester resin, which may then be cured to give a surface coating providing increased strength and resistance to water attack. This coating is burned off either before or during the casting process.\nU.S. Pat. No. 4,361,181 discloses a casting core produced from a granular based substance, in a sugar derivative as binding agent.\nU.S. Pat. No. 3,407,864 discloses use of a fugitive core comprising at least one soluble metal halide salt, and up to about 10% be weight of at least one second constituent selected from the group consisting of borax, magnesium oxide and talc. The patent teaches that it is preferred, for economic reasons, to use sodium or potassium chloride as the metal halide salt. One example of a core made according to the patent contains 95 weight percent sodium chloride, 3 weight percent borax, 1 weight percent magnesium oxide, and 1 weight percent talc.\nU.S. Pat. No. 3,645,491 discloses water-soluble cores for use in the casting of metals, which consist of a water-soluble salt and a synthetic resin. The soluble core may consist of a mixture of the water-soluble salt and a liquid resin binder such as to retain the salt grains in the agglomerate; or alternatively, the soluble core may consist of grains of a water-soluble salt, preferably of a controlled grain size which are coated with a synthetic resin and subsequently dried to a powder. One example of the core composition consists of granular sodium chloride, bonded together in an agglomerate by a synthetic resin of a phenol formaldehyde-furfuryl alcohol type.\nU.S. Pat. No. 3,548,914 discloses soluble cores comprising a mixture of sodium chloride and sodium silicate, or sodium sulfate and sodium silicate.\nU.S. Pat. No. 4,438,804 discloses a water-soluble core prepared from a mixture of sand, potassium carbonate as a first binder, and at least one of barium carbonate and alkali silicate as a second binder.\nAs indicated, the principal problem to which the present invention is addressed is imposed by the use of the high pressure steam and temperatures which occur during the autoclaving step which has been previously discussed. Steam temperatures during autoclaving are typically of the order of 340.degree. to 350.degree. F., accompanied by pressures of about 100 psi. While the prior art water soluble core structures, for example of the types above discussed, have in many instances performed successfully with respect to their ability to be dissolved or disrupted by water, the prior art known cores have not been able to successfully withstand the said autoclaving conditions. Very typically, when a prior art core is thus subjected to autoclaving, the water-soluble salts present in same are found to undergo leaching. This weakens, and indeed may begin to disrupt the entire core structure; and at a minimum, will produce a leached residue of salt on the surface of the core. When such a core is thereupon utilized in the subsequent molding operation, the presence of such salts at the core surface, produces an unacceptable casting.\nAnother difficulty found with certain of the prior art cores, arises from too high thermal expansion characteristics. This can occur, for example, where the sodium chloride content is too high. Lowering the sodium chloride content will reduce the coefficient of thermal expansion, but depending upon the replacement components, leachability may be lost in the fired core--which defeats the purpose of using same.\nIn accordance with the foregoing, it may be regarded as an object of the present invention to provide a nonferrous ceramic core for use in investment casting procedures, which is readily disruptible by water following the casting operation, but which, further, possesses appropriate characteristics as to enable the core to fully withstand the autoclaving conditions which are commonly incident to the investment casting procedures.\nA further object of the invention, is to provide ceramic cores as aforementioned, which have low thermal expansion characteristics, and good leachability as will facilitate removal of the cores following their use in casting."}
{"text": "1. Field of the Invention\nThe present invention relates to the packaging art, and, more specifically, relates to plastic bags or packages formed from extruded polymeric resin sheet material. In particular, the present invention relates to plastic bags or packages of this general type which have zipper closures extruded from polymeric resin material and attached thereto to form a reclosable plastic bag or package.\n2. Description of the Prior Art\nPlastic bags and packages having zipper enclosure extruded from a polymeric resin material, such as polyethylene, are old in the art. Commonly, bags and packages of this type are formed from plastic sheet material on a form-fill-and-seal (FFS) machine and concurrently filled with a consumer item or items. During the formation of the bags or packages, zipper closure strips are introduced between two facing internal portions of the plastic sheet material and sealed thereto.\nRecently, the need to attach zipper closure strips to previously formed or manufactured plastic bags has arisen. The bags in question are formed from a tube of polyethylene. During manufacture, the tube is flattened, and sealed transversely at regularly spaced intervals. The seals ultimately become the bottoms of the bags. Adjacent to each seal, the flattened tube is perforated, so that the bags may be separated from one another. The perforations form the mouths or tops of the bags. The tube, perforated and sealed as described, is wound up in the form of a roll including as many as several thousand individual bags. Individual bags are fed, bottom first, from the roll by automatic bag filling equipment, which sequentially deposits a consumer item or items into each bag and seals each bag transversely below the perforations at the mouth or top thereof.\nWhere bags have previously been formed or manufactured, however, the zipper closure strip or strips must be attached to the outsides of the bags. Such attachment is carried out by some form of heat sealing. Unfortunately, because there is no way to keep the two layers of a previously formed bag separated from one another, there is a tendency for them to be sealed to one another during the attachment process.\nThe present invention provides a solution to this problem."}
{"text": "Integrated circuits interface with their surrounding by input/output (IO) pads that can exchange IO signals. The testing of IO signals and especially their alternate current (AC) characterization is complex. The design of accurate tests and testers for measuring the AC characteristics of IO signals can last months and even years. Setup/hold/skew time between different I/O signals is not simple for measuring on the tester and requires complicated functional tests with strict synchronization solutions.\nThere is a need to provide a tester and method for testing integrated circuits."}
{"text": "The present invention relates generally to method of disinfecting medical materials and more particularly to a method and apparatus for disinfecting medical materials by exposing the materials to radio-frequency waves. The term medical materials encompasses medical waste, veterinary waste, and medical products.\nThe problems with current medical waste handling methods, like the problems of solid waste disposal in general, are becoming increasingly acute. Solid waste is primarily disposed of by burning or by burial in landfill. Both of the methods have severe disadvantages. Burning of solid waste liberates waste particles and fumes which contribute to acid rain and other pollution of the atmosphere. Burying the waste results in possible leaks of toxic chemicals into the surrounding earth and contamination of ground water supplies. Although increasing amounts of solid waste are being recycled, which alleviates the problems of incineration and burial, presently available recycling methods do not provide a complete solution to the disposal problem.\nWaste disposal is of even more urgent concern when the waste comprises possibly infectious medical waste. Such infectious medical waste is a by-product of veterinary and medical care. For example, regulated medical waste consists of: (1) cultures and stocks of infectious agents and associated biological materials; (2) pathological wastes; (3) human blood and blood products; (4) contaminated sharps, including needles, syringes, blades, scalpels, and broken glass; (5) animal waste; (6) isolation waste, including gloves and other disposable products used in the care of patients with serious infections; and (7) unused sharps. These wastes can generally be divided between (a) general medical waste, including cultures and stocks of infectious agents, associated biologicals, pathological waste, and human blood and blood products; (b) veterinary waste, including animal waste; and (c) waste that is predominately plastic, such as the contaminated and unused sharps and isolation waste. The predominately plastic waste also includes metal as well. Hospitals typically segregate waste by types. Contaminated sharps and isolation waste, however, are of special concern as they may carry highly dangerous pathogens such as AIDS virus or hepatitis virus. Sharps in particular have caused widespread public concern when observed washed up on beaches or in public areas.\nHospitals and other generators of medical and veterinary waste employ three methods of waste handling: (a) on-site incineration of the waste, (b) on-site steam autoclaving of the waste followed by later shipment to a landfill for burying, and (c) collection of the waste by a licensed waste hauler with no on-site processing.\nMany hospital incinerators, even those located predominately in urban areas, emit pollutants at a relatively high rate. The Environmental Protection Agency has identified harmful substances in the emissions of such hospital incinerators. They include metals such as arsenic, cadmium and lead, organic compounds, such as ethylene, dioxins and furans, acid gases and carbon monoxide as well as soot, viruses and pathogens. Emissions from these incinerators may be a more significant public health hazard than improper dumping [Steven K. Hall, \"Infectious Waste Management: A Multifaceted Problem,\" Pollution Engineering, 74-78 (August 1989 )].\nAlthough steam autoclaving may be used to sterilize waste before further processing, it is expensive and time consuming. Heat denatures the proteins and microorganisms causing protein inactivation and cell death in a short time. Temperature monitoring devices such as thermocouples, and biological indicators, such as heat resistant Bacillus stearothermophilus spores, may be used to assure effective sterilization.\nU.S. Pat. No. 2,731,208 to Dodd teaches a steam sterilizing apparatus for disposing of contaminated waste which incorporates shredding the waste (\"including paper containers such as used sputum cups, \" col. 1, lines 28-29). Dodd teaches blowing steam into a container full of waste and processing only limited types of items. The Dodd system has the disadvantage of depositing the shredded final mixture into a sewer, which would cause further environmental problems.\nWhether or not the hospital first autoclaves its medical wastes, including broken needles and glass, the waste is then turned over to a licensed waste hauler for transport to a landfill or other depository. U.S. Pat. No. 3,958,936 to Knight discloses compaction of hospital waste for more efficient landfill disposal. Specifically, the reference teaches the application of heat in the range of about 204.degree. C. to 316.degree. C. to hospital and other waste to melt the plastic and convert it into a hard compact block for safer disposal in landfills. The waste is disinfected by the high temperatures, and sharps, such as needles, become embedded in the plastic where they are a reduced mechanical hazard. However, this method suffers from the disadvantage of requiring relatively high temperatures necessitating large energy expenditures and landfill disposal. Metropolitan landfills are becoming filled, and unauthorized dumping is a problem.\nA further area of concern is the sterilization of medical products prior to use. By medical product is meant any product which must be sterilized prior to use in health care. This is exemplified but not limited to needles, syringes, sutures, bandages, scalpels, gloves, drapes, and other disposal items. Many reusable items also must be provided in sterile form. Widespread current sterilization methods include the use of autoclaving, ethylene oxide, and ionizing radiation such as gamma radiation. The heat and humidity of autoclaving are quite damaging to many disposable metal products. Ethylene oxide and ionizing radiation are preferred commercially in those cases.\nIn order to sterilize medical products, poisonous ethylene oxide gas may be used in a closed chamber containing the products to be sterilized. For effective sterilization, not only must the ethylene oxide concentration be controlled carefully, but the temperature, humidity, and porosity of the sterilizer load also must be carefully regulated. Ethylene oxide is relatively slow to dissipate from plastics and its use may require that medical products be stored until the ethylene oxide concentration decreases to a safe level. Ethylene oxide also must be carefully vented to the atmosphere subsequent to the sterilization cycle in order to avoid poisoning operators of the sterilization apparatus.\nIonizing radiation, such as gamma radiation, may be used to sterilize medical products within their packaging; however, it must be administered at such high doses that many plastics become yellow and brittle due to the gamma rays having altered the structure of the polymers of which they are made. For example, U.S. Pat. No. 3,940,325 to Hirao teaches methods for adjusting the formulas of plastics for medical syringes to avoid yellowing and cracking due to exposure to sterilizing gamma radiation. Other substances may also be damaged by exposure to gamma radiation. Such ionizing radiation sterilizes because its high energy photons damage and thereby inactivate the DNA of organisms such as bacteria and viruses. As a result of the inactivation of the DNA, cells lose their ability to reproduce and thereby cause infections. On a large scale industrial basis, ionizing radiation, especially gamma radiation from cobalt 60, has been used to sterilize medical products prior to their use in patients. However, the radiation levels necessary to sterilize may also damage the product being sterilized.\nOther methods have been suggested for sterilization of medical products. For instance, U.S. Pat. No. 3,617,178 to Clouston teaches a method of improving sterilization efficiency by increasing hydrostatic pressure. Elevated hydrostatic pressure causes sterilization resistant bacterial spores to germinate, or begin to grow. However, it has no effect on viruses. Bacterial germination, which converts the bacteria from their environmentally resistant spore form, makes the bacteria more sensitive to radiation, so that lower doses may be employed. Clouston further teaches optimizing the hydrostatic pressure effect by adjusting the temperature up to 80.degree. C. According to Clouston, elevated pressure in heated fluid or moist gas is essential to the method. Elevated temperature alone has a negligible effect. Furthermore, the pressure, heat, or moisture treatment taught by Clouston is intended to cause bacterial spores to germinate thereby rendering them more vulnerable to sterilization techniques, not to sterilize or inactivate microorganisms.\nIn contrast, U.S. Pat. Nos. 4,620,908 to Van Duzer and 3,704,089 to Stehlik teach prefreezing injectable proteins and surgical adhesive prior to irradiation with gamma radiation from cobalt 60 for aseptic manufacture of those materials. U.S. Pat. No. 3,602,712 to Mann discloses an apparatus for gamma irradiation and sterilization of sewage and industrial waste.\nBesides gamma radiation, other types of electromagnetic radiation have been considered as potential sterilants in known systems. Microwaves are increasingly being investigated for rapid sterilization of individual medical devices as well as shredded medical waste. Recently, an experiment showed that metallic instruments could be sterilized in only 30 seconds in a microwave oven (New York Times, \"science Watch Microwave Sterilizer is Developed, \" Jun. 20, 1989). That particular method, however, suffers from the drawback that only a few such metallic instruments can be treated at a particular time. It is not particularly applicable for treatment of medical waste in bulk, and in particular for treatment of medical waste which has been bagged.\nUnited Kingdom Patent No. 1 406 789 to Boucher discloses a microwave system for the surface sterilization of reusable laboratory, medical, and dental instruments in a moist atmosphere at a lower temperature than those presently used and in a shorter time. The system is intended to render aseptic reusable instruments for medical use and generates electromagnetic energy having frequencies between 100 megahertz and 23000 megahertz. Boucher emphasizes that \"his invention deals exclusively with surface sterilization\" and that he \"does not intend to cover such special cases\" as \"`in-depth` sterilization\" (page 1, lines 58-67). Boucher teaches that only through a combination of proper humidification with the thermal and nonthermal effects of microwave radiation can reproducible and satisfactory results be obtained with a wide variety of species, including thermoresistant spores\" (page 1, lines 77-83). Boucher teaches the placement of the object to be sterilized in a gas-tight container with a source of water vapor.\nSoviet Union Patent No. 1,123,705 also discloses a method of sterilizing medical instruments for reuse by UHF treatment. For injection needles it discloses a final temperature of 160.degree. C. to 470.degree. C. and for acupuncture needles it discloses a final temperature of 160.degree. C. to 270.degree. C.\nSystems are also known for treatment of disposable medical waste utilizing microwaves. This system first shreds the waste, sprays the shredded waste with water, and passes the wet shredded waste through a microwave chamber designed to raise the temperature of the wet shredded waste to 205.degree. C. to sterilize it. After the sterilization step, the system compresses the sterilized shredded waste and packages it for shipment to landfills or incinerators (The Wall Street Journal, p. B-3, Apr. 10, 1989). One potential problem with this system is that shredding before sterilization could release infectious particles to the environment and may thus spread contagion. Another problem is the ultimate disposal of the waste; it persists in landfills or may pollute the air when incinerated.\nAlso of interest is a method and apparatus for using microwave frequency electromagnetic fields to heat medical waste to disinfect it. \"Medical Waste Treatment By Microwave Technology\", Norcal Solid Waste Systems. The system includes equipment for receiving the medical waste, shredding it into particle sizes of 1 to 1 1/2 inch linear dimension, and applying steam to the shredded waste to increase its moisture content, as well as to inactivate certain of the microorganisms thereon. The waste is then carried to a microwave treatment area where microwave energy heats the waste to 203.degree. C. for a selected amount of time. A holding area may provide heat sealing. The waste is then recirculated to the steaming station where steam is again applied to inactivate further microorganisms which may still be active in the waste which is shredded and disinfected, disposed in a dumpster for placement in a landfill. It may be appreciated, however, that volumetric heating cannot take place in such a microwave system that the waste has to be scattered in a relatively thin layer on a conveyor belt for treatment by the microwave radiation as the microwave radiation does not adequately penetrate the material. In addition, the material is not enclosed so that there is no substantial transfer of moisture from wet materials to dry materials to aid in the heating within the enclosed system.\nU.S. Pat. No. 3,547,577 to Lovercheck discloses a machine for treating garbage by shredding, compressing the shredded garbage into briquettes, and sterilizing the briquettes with gas. After shredding the garbage is separated into magnetic and nonmagnetic portions. The sterilization step employs ethylene oxide gas which requires temperature control. The briquettes are maintained at a temperature of about 54.degree. C.\nFurther, microwaves are limited in their penetration and are ineffective for heating when applied to large scale, boxed medical waste of the type which comprises the waste disposal problem today. Microwaves do not heat very effectively because they do not penetrate very deeply. Most of the heat is generated near the surface and quickly dissipates into the surroundings, in part because it is not well conducted into the center portions of the boxed medical waste. In contrast, radio-frequency waves at relatively low frequency can penetrate boxed medical waste more deeply.\nIt also is known in the art that thermal radiation treatment of bacterial spores and other pathogens may allow greatly reduced ionizing radiation dosage to accomplish sterilization of a given population. For instance, in \"Thermoradiation Inactivation Of Naturally Occurring Bacterial Spores In Soil, \" M. C. Reynolds et al., Applied Microbiology, Vol. 28, No. 3, September 1974, it is disclosed that bacterial spores may be inactivated by heating them with dry heat and exposing them to ionizing radiation from a cobalt 60 source allow greatly reduced treatment times over the use of either dry heat or radiation alone.\nAn attempt to elucidate a model for such behavior is set forth in J. P. Brannen, \"A Kinetic Model For The Biological Effects Of Ionizing Radiation\", Sandia Laboratories, SAND74-0289 (October 1974).\nHeat and radiation inactivation of bacteria are discussed at \"Progress Report Beneficial Uses Program, Period Ending Dec. 31, 1976\", Waste Management and Environmental Programs Department, Sandia Laboratories, SAND77-0426 (1977), where it is taught that viruses in sewage sludge may be destroyed by evaporation. Heat inactivation may be used to destroy Salmonella enterititis ser. montevideo. Streptococcus bacteria may be destroyed by using ionizing radiation at a dose of about 140 kilorads.\nThe use of cesium 137 to inactivate pathogens in sludge is discussed in \"Sludge Or Radiation Disinfection For Beneficial Use\", Applied Biology And Isotope Utilization Division 4535, \"General Description Of The Sludge And Radiation Process\", SAND80-2744 (December 1980), where it is disclosed that cesium-137, emitting gamma radiation may be used to inactivate pathogens in sewage sludge. See also, \"Use Of Cesium-137 To Process Sludge For Further Reduction Of Pathogens, Sludge Or Radiation Disinfection For Beneficial Use,\" Disease Control Requirements For Various Sludge Uses, Applied Biology and Isotope Utilization Division 4535, SAND80-2744 (December 1980), which discloses that in order to render sewage sludge safe, in particular for certain agricultural usages, irradiation may be used as an add-on process in conjunction with sterilization where sludge is maintained at 30 min. at a temperature of at least 70.degree. C. In each of the aforementioned papers, it may be appreciated that the sludge which is being treated is substantially homogeneous in its dielectric characteristics and, thus, in its heating characteristics.\nThe gamma irradiation equipment commonly used and disclosed in this application is of the type disclosed in \"Gamma Processing Equipment\", AECL Industrial and Radiation Division (January 1987).\nThe dual plate about 12 megahertz plate type radio-frequency heater is of the type disclosed in \"Dielectric Heating\", PSC, Inc., which, although undated, constitutes prior art to this application.\nLike microwaves, radio-frequency waves are a form of electromagnetic energy. They also transfer energy directly into materials, primarily by the interaction of their time-varying electric fields with molecules. Radio-frequency waves may be applied by connecting a radio-frequency alternating current to a pair of electrodes. Between the two electrodes an alternating radio-frequency electromagnetic field having a time-varying electric field component is established. When objects are placed between the electrodes in the time-varying electric field, the time-varying electric field partially or completely penetrates the object and heats it.\nHeat is produced when the time-varying electric field accelerates ions and electrons which collide with molecules. Heat also is produced because the time-varying electric field causes molecules, and particularly those with a relatively high electric dipole moment, to rotate back and forth as a result of the torque placed upon them by the time-varying electric field. Most large molecules, or molecules with evenly distributed charge, have relatively low or nonexistent dipole moments and are not very much affected by the radio-frequency time-varying electric field. Small molecules, in particular with polar groups, have relatively large electric dipole moments and thus have relatively large torques exerted upon them by the time-varying electric field. In particular, highly polar molecules, like water, experience relatively large torques and as a result are rotated by the time-varying electric field, thereby transferring mechanical energy to their surroundings as internal energy or heat. Lower frequency time-varying electric fields penetrate deeply and heat objects more evenly. Relatively high frequency time-varying electric fields do not penetrate as deeply, but heat more rapidly the portions of objects they interact with.\nIt should be noted that a time-varying electric field is always accompanied by a time-varying magnetic field, except where destructive cancellation occurs with interference patterns. For most materials being considered here, the principal heating mechanism arises from the electric fields. These fields can cause both ohmic heating via induced ionic currents and dielectric heating via molecular stressing from the internal electric fields. For very moist materials, the presence of the accompanying time-varying magnetic field can also induce eddy-currents which can also heat the material. Also, some type of combined effect of magnetic fields and heat may occur. While the ensuing discussion is presented in context of an electric field effect, it should be understood that the effects of accompanying time-varying magnetic field are defined here for simplification as part of the electric field phenomena.\nBecause different materials are composed of different types of molecules with differing electric dipoles, they heat at different rates when exposed to a given time-varying electric field. For example, plastics, which are formed of very large polymer molecules, are not heated by time-varying electric fields as rapidly as water. Metal objects may or may not be easily heated when exposed to time varying electric fields either in the radio-frequency or microwave region. The high conductivity of the metal objects tends to short cut the electric fields and rescatter them. As a consequence, there are many conditions where metal objects are difficult to heat, as exemplified by the metal liner of the interior microwave ovens. On the other hand, such time-varying fields can also induce substantial currents which flow on the outside of the metal objects. Under certain circumstances heating effects will occur on the surface of the metal object which, in the case of a small needle, the heat is readily diffused to the interior. In addition, the presence of long, thin metal objects in an electric field causes enhancement of the electric field intensity near the ends of the metal objects and a diminution or shadowing of the fields near the middle. Thus, if the electric field is parallel to the axis of the metal object, strong electric fields will exist near the tips and weak electric fields will exist near the center of the rod or needle. Such field enhancements can lead to arcing and possible fires. In addition, the field suppression or shadowing of such metal objects is also an unwanted feature if the presence of a single electric field vector is relied upon in its entirety to provide the sterilization. The failure of the radio-frequency electromagnetic field to penetrate the object causing surface heating only, or the opposite failure of the materials to absorb the electric field energy, causes uneven heating of the medical waste. The uneven heating is exacerbated because the medical waste usually comprises mixed materials which are difficult to heat effectively using radio-frequency energy due to the presence of areas of high field absorption, such as are due to metals and concomitant shadowing and cold spots. In addition, similar but less pronounced absorption effects are found with water molecules. Thus, when heterogeneous or mixed medical wastes have wet and dry portions, it may be seen that only the wet portions of such material would be heated. Mixed loads such as hospital wastes were considered impossible to disinfect by radio-frequency energy because the waste contains a wide variety of materials, each having different dielectric properties. A great concern was that the presence of a sufficient number of metallic sharps would lead to arcing, causing ignition of the accompanying dry wastes. Another concern was that even if fire was not started, the differential energy absorption of fluids and sharps would leave dry objects undisinfected.\nIn fact, other attempts to kill microorganisms with radio-frequency energy have been considered unsuccessful. In his 1980 review, \"Effects Of Microwave Irradiation On Microorganisms\", Advances in Applied Microbiology 26:129-45, Chipley cites an experiment of applying radio-frequency energy to bacteria and viruses which grow on tobacco. The experiment found no effect of the radio-frequency energy on the bacteria and viruses. In another study of radio-frequency energy on contaminated liquid food, there was no showing of \"selective killing effect\" except when ethanol was added.\nIn the same review, Chipley cited numerous tests of microwaves on microorganisms and concluded that \"results of tests for viability of B. subtilis spores also showed identical death curves compared with those obtained by conventional heat.\" On the other hand, however, Chipley cites several references which support the view that microwave irradiation has collateral thermal and nonthermal effects. [For example, Culkin and Fung (1975) found that microbial destruction occurred at reduced temperatures and shorter time periods when the material was exposed to microwaves as compared to conventional heating methods. Wayland et al., 1977 also demonstrated the interdependence of heat and microwaves effects in the studies of spores of B. subtills.\nU.S. Pat. No. 2,114,345 to Hayford discloses a radio-frequency applicator with electroscopic control for destroying bacteria in bottled beer and similar products. Hayford teaches an apparatus for sterilizing a series of small objects. The radio-frequency field must be constantly readjusted by the electroscopic control. There is no teaching or suggestion that large scale sterilization of heterogeneous waste could be carried out.\nU.S. Pat. No. 3,948,601 to Fraser et al. teaches the indirect use of radio-frequency energy in sterilizing medical and hospital equipment as well as human waste. The reference teaches the use of radio-frequency energy for heating gases, particularly argon, and exciting them so that they ionize into a plasma having a temperature of approximately 100.degree. C. to 500.degree. C. The reference teaches that a cool plasma at a temperature of only 25.degree. C. to 50.degree. C. and very low pressure may effectively sterilize an article. However, sterilization by plasma does not suggest the direct use of radio-frequency waves in sterilization since it is the chemical reactive effect of the plasma which presumably performs the sterilization function rather than the direct or thermal effects of radio-frequency energy on pathogens contained on the material. It may be appreciated that only those portions of the equipment and waste actually contacted by the plasma would be treated.\nReprocessing of the waste, and especially medical waste, is also vital for several reasons. Even if the medical waste has been rendered harmless or innocuous by the destruction of any pathogens associated therewith, there is still the problem of the disposal of the bulk material including the plastics, the sharps, and fibrous material such as gowns, diapers, and the like. The material is relatively bulky and landfills, particularly in many urban areas, have become filled. In addition, older landfills may leak and nonpathogenic but chemically polluting substances may leak into surrounding ground water, causing health hazards. Thus, burying the sterilized medical waste is becoming less attractive. Further, merely burning the sterilized medical waste can pollute the atmosphere and cause acid rain. Current reprocessing technology should be employed to process the disinfected medical waste for effective utilization and proper disposal. What is needed is a method for disinfecting the medical waste and destroying the pathogens thereon and disposing of the disinfected waste in a manner which is harmless to health care workers, waste handlers, and the public at large.\nA series of investigations has been undertaken as to sterilization, especially for food. This has resulted in patents or inventions wherein the material to be treated is housed in a microwave transparent container such that the material can be heated at vapor pressures which coexist with temperatures of 120.degree. C. These include Gray U.S. Pat. No. 3,494,723; Nakagawa U.S. Pat. No. 4,808,782; Stenstron U.S. Pat. No. 4,808,783; Landy U.S. Pat. No. 3,215,539; Utosomi U.S. Pat. No. 3,885,915; and Fritz U.S. Pat. No. 4,775,770. All of these patents disclose heating homogeneous material in some form of pouch or pressure container where the material, typically food, is homogeneous. They do not address the special problem considered here where the material is heterogeneous and contains sharps, moist materials and dry materials."}
{"text": "The present invention relates to an information control method for a securities dealing support system such as a trading support system which supports, or assists in the dealing of securities such as bonds and debentures, and to a securities dealing support system itself.\nA conventional dealing system for bonds, debentures or the like is disclosed in Japanese Patent Laid open Publication JP-A-61-86867. With this system, customers who intend to deal with bonds or debentures are supplied with a plurality of type of data made at a host computer in menu form, or with the information requested by customers. The system automatically processes securities dealing proposed by customers. Information exchange between investors and dealers and/or traders is carried out in the following manner. The term \"dealers and/or traders\" is represented by \"dealer/trader\" hereinunder where applicable.\nThere is shown in FIG. 2 the flow of information about the decision support for the dealer/trader who deals with bonds and debentures, and about the dealings between the dealer/trader and customers. Namely, in accordance with the market information at a stock exchange 21 or the like dealing with bonds, debentures and the like, the dealer/trader 23 consults customers 24 for securities dealing. In accordance with the securities dealing information supplied from customers 24, the dealer/trader 23 sends an order to the market. The information has been transferred heretofore via a telephone between dealers and traders or between traders and customers. The dealing information has been supplied to the dealer/trader in such a manner that a computer calculates the dealing support information in real time on the basis of market information and outputs the information to the display 22 of a work station at the dealer/trader. There is shown in FIGS. 3A and 3B an example of a conventional trading support computer system. In the system shown in FIG. 3A, market information is received at a computer 31 which selects desired information and sends it onto a Local Area Network (LAN) 32. A terminal device 33 picks up desired information from the information flowing on LAN 32 to thereby supply necessary information for dealing with bonds, debentures or the like. The dealer/trader communicates with customers via a communication system 34 while observing the screen of the terminal device 33.\nIn the system shown in FIG. 3B, A computer 31 receives market information and calculates desired securities dealing support information which is supplied to a terminal device 33. The dealer/trader communicates with customers via a communication system 34 while observing the screen of the terminal device 33."}
{"text": "The present invention concerns a base station in a mobile telecommunication access network such as the UMTS terrestrial radio access network (UTRAN). The UTRAN is illustrated in FIG. 1 and comprises at least one Radio Network System 100 connected to the Core Network (CN) 200. The CN is connectable to other networks such as the Internet, other mobile networks e.g. GSM systems and fixed telephony networks. The RNS 100 comprises at least one Radio Network Controller 110. Furthermore, the respective RNC 110 controls a plurality of Node-Bs 120,130 that are connected to the RNC by means of the Iub interface 140. Each Node B, also referred to as base station, covers one or more cells and is arranged to serve the User Equipment (UE) 300 within said cell. Finally, the UE 300, also referred to as mobile terminal, is connected to one or more Node Bs over the Wideband Code Division Multiple Access (WCDMA) based radio interface 150.\n3GPP Release 6 has recently been updated with the “Enhanced Uplink” concept, which includes a new uplink transport channel, E-DCH. A new scheduler located in the Node B that performs transmission resource management based on the utilisation of the allocated radio resources is introduced as an alternative or complement to the existing packet scheduler located in the RNC. Combined with fast L1 HARQ schemes, the proposed algorithm can provide a cell throughput gain. Since the packet scheduler functionality is located in the Node B, a fast scheduling concept is introduced. Fast scheduling from the Node B denotes the possibility for the Node B to control when a UE is transmitting and at what data rate. The Node Bs can assign scheduling grants to the UEs, where these grants are based on both the transmission resource availability and the requested need for transmission resources, i.e. the scheduling requests from the UEs.\nHARQ is a more advanced form of an ARQ retransmission scheme. In conventional ARQ schemes the receiver checks if a packet is received correctly. If it is not received correctly, the erroneous packet is discarded and a retransmission is requested. With HARQ the erroneous packet is not discarded. Instead the packet is kept and combined with a result of the retransmission. That implies that even if both the first transmission and the retransmission are erroneous, they may be combined to a correct packet. This means that fewer retransmissions are required.\nOn E-DCH the Node B scheduler has no direct information about the data that to be transmitted from the UEs. Thus the UEs are required to indicate the amount of data available, the priority of the data, the transmitter power available etc. to the Node B through scheduling requests. When the Node B has received the scheduling request from the UE and has decided to schedule the UE based on the received scheduling requests, it may transmit a grant, also denoted scheduling grant indicator herein, to the UE, indicating the amount of data or actually with which power the UE is allowed to transmit.\nA particular aspect of relevance for the present invention is the fact that enhanced uplink, i.e. the E-DCH supports soft-handover. Soft-handover implies that a UE is connected to multiple base stations simultaneously. Thus, a UE in soft handover is power-controlled from multiple cells, and supported by data reception at multiple cells (i.e. macro diversity is utilized). Power control from multiple cells is needed to limit the inter-cell interference, while macro-diversity gains can be achieved by receiving data at multiple cells. FIG. 2 illustrates a scenario when a UE is in soft handover in a UMTS network as shown in FIG. 1. The network comprises base stations connected to a RNC 208, wherein the RNC 208 is further connected to a CN 210. The UE 202 is connected to the base stations 204 and 206 simultaneously.\nIn soft-handover in enhanced uplink, one cell is selected by the RNC to act as the Serving Cell, and the Node B in control of the Serving Cell is here denoted the Serving E-DCH Node B. The cells connected to the same UE is referred to as the E-DCH active set. The Serving E-DCH Node B controls the transmission resources of the UE, i.e. it is allowed to grant requested transmission resources. The Non-serving E-DCH Node Bs (i.e. Node Bs not in control of the serving cell) are not able to monitor the grants given by the E-DCH serving Node B and has thus no knowledge of the granted transmission resources. The Serving E-DCH Node B should typically be the Node B with the strongest uplink from the UE, but the E-DCH serving Node B may also be chosen differently. One alternative is to tie the E-DCH serving cell to the downlink HS-DSCH serving cell. Tying the E-DCH serving cell selection to the HS-DSCH serving cell may increase the likelihood that the strongest uplink in the active set is governed by a non-serving E-DCH Node B.\nThe present invention considers problems arising from such configuration that may result in that the uplink of the serving Node B is not the strongest link in the active set, i.e. another Node B of the active set receives the strongest signal from the UE in soft handover.\nThe E-DCH scheduling is mainly controlled from the serving E-DCH Node B, which can assign Absolute Grants (AG) to the UE's. These Absolute Grants limit the maximum transmission resources, e.g. power, the UE is allowed to use. Within this restriction, the final selection of data-rate is then performed by the UE itself, based on the data available in its buffers and on the available UE power. Alternatively, the serving E-DCH Node B can use Relative Grants to control the transmission rates of the UEs. The Relative Grants from the serving E-DCH cell can take three values: Up, Hold, and Down. However, the non serving Node Bs do not receive any information of limitations indicated by the absolute grant or relative grant from the serving E-DCH cell and is therefore not aware of the future processing resource need of the UE.\nThe scheduling control from non-serving E-DCH Node Bs is mainly intended for inter-cell interference suppression and system stability control. The Relative Grants that can be sent from the non-serving E-DCH Node Bs therefore take only two values: Hold and Down. With these relative grants, a Node B can reduce the interference contribution from UEs, which are not primarily controlled from this Node B.\nThus, the present invention deals with the processing resource allocation problem of the Node B that depends on that the non-serving Node Bs are not aware of the amount of transmission resources that is granted from the Node B controlling the serving cell.\nFrom a Node B internal hardware allocation point of view, there is a significant difference between the serving E-DCH Node B and the non-serving Node Bs. The serving E-DCH Node B has information about the scheduling grant sent to the UE and therefore knowledge about the maximum amount of hardware resources needed for processing transmissions from this particular UE. However, a Node B that is not in control of the Serving E-DCH cell has much more limited means to predict the processing resources needed for the UE in question. This resource allocation problem is further complicated by the fact that the Node B processing resource allocation typically takes some time (e.g. 10-50 ms), meaning that a predictive processing resource allocation would be necessary to ensure that the reception of a transmission from a UE can be successfully completed.\nBelow are two existing solutions to this problem described:\nAccording to a first example, processing resources are over-allocated and rate-limitation of terminals in soft handover may also occur. In this brute-force solution, sufficient Node B processing resources for the highest possible data rate are always allocated from the non-serving Node B. To reduce the need of hardware resources, the maximum bit rate of UEs in soft-handover may have to be limited.\nAccording to a second example, the non-serving Node B could under-allocate processing resources, knowing that it may not be able to decode the first few Transmission Time Intervals (TTIs) of a UE transmission, in case the UE starts at a rate higher than estimated. Once the UE starts to transmit at a high data rate, the non-serving Node B can reallocate processing resources to this UE, assuming that it will continue to transmit for some time. Non-serving Node Bs may also try to listen to the scheduling requests from the UE to the serving cell to get some information about the amount of transmission resources the UE may need.\nBoth approaches are, however, equipped with obvious drawbacks:\nThe solution according to the first alternative may result in low utilization of available hardware resources, costly deployments with a necessity to deploy large pools of Node B processing resources and/or tight bit-rate restrictions for users in soft handover. This is neither desirable nor acceptable.\nThe opportunistic approach in the solution according to the second example may result in loss of the macro-diversity gain, since the first few HARQ transmission attempts at a non-serving Node B are lost. In addition, the second solution may interact with Outer-Loop Power-Control, which typically is operated on the number of HARQ transmissions. An increase in the number of HARQ transmissions due to hardware limitations may result in an unnecessary increase in the uplink SIR target with a capacity loss as a consequence."}
{"text": "1. Field of the Invention\nThe invention relates to thermal insulation, especially for high-temperature batteries, wherein the thermal insulation has at least one hermetically sealed cavity which contains insulating material.\n2. Description of the Prior Art\nSuch thermal insulation finds use in preventing thermal losses of devices which are used in energy engineering. Among others, they are intended for high-temperature batteries such as batteries of the alkali metal and chalcogen type which are usually surrounded by thermal insulation in order to prevent the storage cells from cooling down, particularly in the operating intervals. Up to now, insulations used for such batteries are built, for instance, of glass or mineral wood. To obtain a sufficiently large effect, considerable wall thicknesses must be provided in these insulations, and specifically in those cases where the storage battery operates at elevated temperatures, for example at 350.degree. C., and where these temperatures are to be maintained during the operating intervals of the battery, which may be hours. Since such thick-walled insulation increases the dimensions and/or the weight of the storage battery considerably, the energy storage density, i.e. the electric energy that can be stored per weight of volume unit, is small. This is a disadvantage, particularly for electrochemical batteries which are to be used as the power supply of electrically operated vehicles.\nAnother thermal insulation for such electrochemical storage cells is known. It is formed essentially by an evacuated cavity, within which metal foils are arranged. These metal foils are made particularly of aluminum and are arranged at a predeterminable distance from each other. The cavity is formed by metallic walls with a small thermal expansion coefficient. In order to prevent bending of the boundary walls (due to the existing vacuum) inward, bar-shaped supports between parallel boundary walls are arranged in the interior of the cavity. These support elements, however, cause a large heat flow from the inside of the insulation toward the outside. Thereby, the low thermal conductivity of the insulation which is obtained by the metallic foils, is largely lost again."}
{"text": "1. Field of the Invention\nThe present invention relates to a cordless power transmission system, power transmission terminal and electrical appliance, which transfer power by electromagnetic waves.\n2. Description of Related Art\nAs a conventional cordless power transmission system, there is known an LSI that is used for non-contact card or the like. It receives very feeble microwaves emitted into a communication space within a certain range, and operates by converting them into electric power. Such a system is applied to simple data transfer in ticket inspection at a ski lift, or in an intelligent transport system, for example.\nSince the conventional cordless power transmission system supposes a single device within a narrow range, it cannot designate a specific device when there are multiple LSI cards within the range of the very feeble microwaves. Thus, the system can bring about erroneous operations in such a situation, presenting a problem of being unsuitable for large scale data communications or radio power transmission. Besides, the system is lacking in versatility with only limited applications.\nThe present invention is implemented to solve the foregoing problems. It is therefore an object of the present invention to provide a cordless power transmission system, power transmission terminal and electrical appliance suitable for large scale data communications and radio power transmission, and capable of transmitting and receiving electric power stably without any malfunctions even from separate places, thereby making all electrical appliances in a room wireless, for example.\nAccording to a first aspect of the present invention, there is provided a cordless power transmission system comprising: at least one power transmission terminal connected to a network; at least one electrical appliance supplied with power from the power transmission terminal by electromagnetic waves; a network host connected to the network for managing the electrical appliance via the power transmission terminal; a default power transmitter installed in the power transmission terminal for transmitting default power to the electrical appliance by diffractive electromagnetic waves; a first data transceiver installed in the electrical appliance for receiving default power, and for transmitting unique data and required power intensity data to the power transmission terminal by diffractive electromagnetic waves; a second data transceiver installed in the power transmission terminal for receiving the unique data and required power intensity data from the electrical appliance, and for transmitting them to the network host; and a power transmitter installed in the power transmission terminal for transmitting ID data to be assigned to the electrical appliance and power with a required intensity to the electrical appliance by rectilinear electromagnetic waves, wherein the network host obtains the ID data to be assigned to the electrical appliance from the unique data of the electrical appliance, and obtains the required intensity of the power to be transmitted from the power transmission terminal from the required power intensity data sent from the electrical appliance.\nHere, the first data transceiver may continuously transmit the ID data, which is assigned to the electrical appliance by the network host, to the power transmission terminal by the diffractive electromagnetic waves while the electrical appliance is in a normal operation mode, and the second data transceiver may notify the network host of the ID data received by the power transmission terminal.\nThe cordless power transmission system may further comprise a plurality of power transmission terminals, each of which includes a data collector for collecting intensity data of the electromagnetic waves arriving at the power transmission terminal from the electrical appliance, wherein the network host may obtain a 3-D position of the electrical appliance from the intensity data of the electromagnetic waves collected by the power transmission terminals, and assigns a part of the required intensity power to each of the power transmission terminals.\nThe cordless power transmission system may comprise at least two power transmission terminals, each of which includes a reflectance detector for detecting reflectance of the rectilinear electromagnetic waves while the power transmission terminals are transmitting required power by the rectilinear electromagnetic waves, and the second data transceiver of each of the power transmission terminals may notify the network host of the reflectance, wherein the network host has the power transmission terminals other than the power transmission terminal whose reflectance of the electromagnetic waves varies substantially, transmit the required power.\nThe network host may be installed in each room, and the network hosts may be interconnected by a cable to exchange data between the rooms.\nThe network hosts may be connected to an external communication circuit via a fire wall.\nThe diffractive electromagnetic waves may consist of radio waves; and the rectilinear electromagnetic waves may consist of light waves.\nThe power transmitter of the power transmission terminal may convert AC power into light power with the required intensity, transmit it to the electrical appliance by light waves, and carry out information communication with the electrical appliance by the light waves; the default power transmitter may convert the AC power into radio power, and transmit the default power to the electrical appliance by the radio waves; and the second data transceiver may carry out information communication with the electrical appliance by the radio waves.\nThe electrical appliance may further comprise: a light section for converting light power transmitted from the power transmission terminal into electric power, and for carrying out information communication with the power transmission terminal by light waves, wherein the first data transceiver may convert the radio power transmitted from the power transmission terminal into electric power, and carry out information communication with the power transmission terminal by the radio waves.\nAccording to a second aspect of the present invention, there is provided a power transmission terminal that supplies power to an electrical appliance by electromagnetic waves in accordance with a command sent from a network host connected to a network, the power transmission terminal comprising: a default power transmitter for supplying the electrical appliance by diffractive electromagnetic waves; a data transceiver for receiving unique data and required power intensity data associated with the electrical appliance from the electrical appliance by diffractive electromagnetic waves, and notifies these data to the network host; and a power transmitter for transmitting ID data and power with a required intensity to the electrical appliance by rectilinear electromagnetic waves, the ID data and the required intensity of the power being assigned to the electrical appliance and to the power transmission terminal, respectively, by the network host.\nAccording to a third aspect of the present invention, there is provided an electrical appliance that receives power by electromagnetic waves transmitted from a power transmission terminal connected to a network host through a network, the electrical appliance comprising: a default power receiver for receiving default power from the power transmission terminal by diffractive electromagnetic waves; a data transmitter for transmitting unique data and required power intensity data associated with the electrical appliance to the power transmission terminal by diffractive electromagnetic waves; and a power receiver for receiving ID data and power with a required intensity from the power transmission terminal by rectilinear electromagnetic waves, the ID data and the required intensity of the power being assigned to the electrical appliance and to the power transmission terminal, respectively, by the network host."}
{"text": "Carburetors are devices that can be used to mix fuel and air to power combustion engines typically including gasoline powered internal combustion spark ignited engines. A carburetor may include a fuel metering system that helps to control the amount of fuel supplied to air flowing through a mixing passage or main bore of the carburetor for mixing the fuel with air and supplying the mixture to the engine. Some metering systems employ a diaphragm that oscillates or reciprocates during operation to open and close a metering valve admitting fuel to a chamber from which it is supplied to the passage for mixing with air. In use, the large number of cycles experienced by such a diaphragm which typically physically interacts with other components of the metering system such as a valve actuating lever, and continuous exposure to solvent containing fuels, can result in a harsh operating environment that causes wear, degradation and ultimately failure of the diaphragm. In a gasoline powered spark ignited internal combustion so-called small engine the diaphragm must fully open the valve when subjected to only a small pressure differential which is typically a maximum negative pressure of −0.9956 kPa or −0.1444 pounds per square inch and usually about −0.50 kPa or −0.0725 pounds per square inch (psi). This very small differential operating pressure also requires that the portion of the diaphragm within a fuel metering chamber be very flexible particularly since such diaphragm may have a surface area within the metering chamber in the range of about 0.5 square inch to 1.0 square inch."}
{"text": "In use in a gas turbine engine, airfoil components, such as gas turbine engine blades and vanes, are operated at extremely elevated temperatures and as a result are internally cooled by ducted airflow through one or more cooling passages located internally in the airfoil. Such components typically are externally coated with a protective coating that is resistant to high temperature degradation. The walls of the internal passages of such components also typically can be protectively coated to this same end. Over time, such airfoil components exhibit coating wear, cracking, corrosion, and other degradation to the coating and/or airfoil substrate that necessitates repair or replacement of the component. Repair, rather than replacement, is the preferred procedure for extending the service life of such components from an economic standpoint due to their high initial cost.\nIn typical repair procedures, engine-run airfoil components are subjected to external and possible internal recoating with protective coating materials using well known high temperature coating techniques, such as pack, out-of-pack, chemical vapor deposition, or plasma spray coating to form protective external and internal coatings. It is important prior to such coating operations that any deposits of foreign material accumulated or built-up on the external and/or internal surfaces of such airfoil components be removed to avoid harmful metallurgical contamination of the repaired airfoil component. Dirt build-up reduces airflow cooling which can create a \"hot spot\" in the airfoil, possibly resulting in blade failure. Removal of deposits from external surfaces of engine-run airfoil components is readily achieved as a result of the accessibility of the external surfaces to chemical cleaning agents, such as caustic solutions, or to grit blasting agents. However, removal of deposits from internal passages of engine-run airfoil components is rendered difficult by virtue of their small size and oftentimes convoluted nature that can hinder access of cleaning solution throughout the passage.\nOne technique developed to remove difficult-to-access internal passages of engine-run airfoil components employs a high pressure autoclave procedure wherein the engine-run components to be cleaned are disposed inside an autoclave containing a caustic cleaning solution, such as an aqueous 45% KOH solution. The autoclave is heated to elevated temperature greater than 400.degree. F. and pressurized to 200 psi for a prolonged time (e.g. 8-24 hours) to remove deposits from the internal airfoil passages. However, this autoclave procedure may not be effective to adequately remove heavy deposits from the internal airfoil passages such that the components are scrapped rather than subjected to remaining repair operations that would enable reuse of the repaired component in gas turbine engine service."}
{"text": "Streptococcus pneumoniae (S. pneumoniae) is a common, spherical, gram-positive bacterium. Worldwide it is a leading cause of illness among children, the elderly, and individuals with debilitating medical conditions (Breiman, R. F. et al., 1994, JAMA 271: 1831). S. pneumoniae is estimated to be the causal agent in 3,000 cases of meningitis, 50,000 cases of bacteremia, 500,000 cases of pneumonia, and 7,000,000 cases of otitis media annually in the United States alone (Reichler, M. R. et al., 1992, J. Infect. Dis. 166: 1346; Stool, S. E. and Field, M. J., 1989 Pediatr. Infect. Dis J. 8: S11). In the United States alone, 40,000 deaths result annually from S. pneumoniae infections (Williams, W. W. et al., 1988 Ann. Intern. Med. 108: 616) with a death rate approaching 30% from bacteremia (Butler, J. C. et al., 1993, JAMA 270: 1826). Pneumococcal pneumonia is a serious problem among the elderly of industrialized nations (Käyhty, H. and Eskola, J., 1996 Emerg. Infect. Dis. 2: 289) and is a leading cause of death among children in developing nations (Käyhty, H. and Eskola, J., 1996 Emerg. Infect. Dis. 2: 289; Stansfield, S. K., 1987 Pediatr. Infect. Dis. 6: 622).\nVaccines against S. pneumoniae have been available for a number of years. There are a large number of serotypes based on the polysaccharide capsule (van Dam, J. E., Fleer, A., and Snippe, H., 1990 Antonie van Leeuwenhoek 58: 1) although only a fraction of the serotypes seem to be associated with infections (Martin, D. R. and Brett, M. S., 1996 N. Z. Med. J. 109: 288). A multivalent vaccine against capsular polysaccharides of 23 serotypes (Smart, L. E., Dougall, A. J. and Gridwood, R. W., 1987 J. Infect. 14: 209) has provided protection for some groups but not for several groups at risk for pneumococcal infections, such as infants and the elderly (Mäkel, P. H. et al., 1980 Lancet 2: 547; Sankilampi, U., 1996 J. Infect. Dis. 173: 387). Conjugated pneumococcal capsular polysaccharide vaccines have somewhat improved efficacy, but are costly and, therefore, are not likely to be in widespread use (Käyhty, H. and Eskola, J., 1996 Emerg. Infect. Dis. 2: 289).\nAt one time, S. pneumoniae strains were uniformly susceptible to penicillin. The report of a penicillin-resistant strain of (Hansman, D. and Bullen, M. M., 1967 Lancet 1: 264) was followed rapidly by many reports indicating the worldwide emergence of penicillin-resistant and penicillin non-susceptible strains (Klugman, K. P., 1990 Clin. Microbiol. Rev. 3: 171). S. pneumoniae strains which are resistant to multiple antibiotics (including penicillin) have also been observed recently within the United States (Welby, P. L., 1994 Pediatr. Infect. Dis. J. 13: 281; Ducin, J. S. et al., 1995 Pediatr. Infect. Dis. J. 14: 745; Butler, J. C., 1996 J. Infect. Dis. 174: 986) as well as internationally (Boswell, T. C. et al., 1996; J. Infect. 33: 17; Catchpole, C., Fraise, A., and Wise, R., 1996 Microb. Drug Resist. 2: 431; Tarasi, A. et al., 1997 Microb. Drug Resist. 3: 105).\nA high incidence of morbidity is associated with invasive S. pneumoniae infections (Williams, W. W. et al., 1988 Ann. Intern. Med. 108: 616). Because of the incomplete effectiveness of currently available vaccines and antibiotics, the identification of new targets for antimicrobial therapies, including, but not limited to, the design of vaccines and antibiotics, which may help prevent infection or that may be useful in fighting existing infections, is highly desirable."}
{"text": "1. Field of the Invention\nThe present invention relates to a color prediction method, a color conversion lookup table generation method, an image processing method and an image processing apparatus, and more particularly, to color correction technology for image forming data in color image formation.\n2. Description of the Related Art\nIn general, in an inkjet recording apparatus or an image forming apparatus based on an electrophotographic method, a plurality of coloring materials, such as yellow (Y), magenta (M), cyan (C) and black (K) are used. The color space achieved by these coloring materials is different to that of the original image data, and therefore color conversion (color correction) processing is carried out in order to faithfully reproduce the colors and tones represented by the original image data.\nOn the other hand, depending on the characteristics of the output medium for image formation, there are limits on the amount of coloring material that can be received by the medium. For example, when forming a color image on a paper medium using an inkjet recording apparatus, deformation of the paper medium occurs if ink of an amount exceeding the amount that can be absorbed by the paper medium is deposited thereon. Due to this deformation of the paper medium, there is a risk of the occurrence of conveyance abnormalities of the paper medium, and decline in the quality of an image formed on the paper medium.\nFurthermore, even if the paper does not deform, when an excessive amount of coloring material is used, coloring material may not be fixed satisfactorily on the medium, leading to degradation of the film properties and reduced resistance to rubbing, and problems such as transfer to the rear surface of other media may occur when media are stacked after image formation.\nJapanese Patent Application Publication No. 2004-58623 discloses a method according to which, when generating a color separation table (a table used for color separation processing to calculate a coloring material use amount from signal values) on the basis of data obtained by color measurement of patches, if there is a patch having a coloring material use amount greater than a maximum overall coloring material use amount specified on the basis of the ink absorption characteristics of a recording medium, and the like, then a color separation table within the range of the maximum overall coloring material use amount is generated by compressing the signal values in such a manner that all of the patches come within the range of the maximum overall coloring material use amount."}
{"text": "The present invention relates to a transfer apparatus having a primary part, at least one first or one second secondary part, which are each magnetically coupled to the primary part for drive purposes, and a regulation device for regulation of a drive of the secondary parts. The present invention also relates to a corresponding method for magnetic driving of at least one first and one second secondary part with a primary part.\nThe following discussion of related art is provided to assist the reader in understanding the advantages of the invention, and is not to be construed as an admission that this related art is prior art to this invention.\nA transfer apparatus is used, for example, for transportation and precise positioning of workpiece carriers with workpieces arranged on them. A transfer apparatus such as this can also be used as part of and/or along a conveyor path. One or more processing stations and one or more positioning modules, suitably with a positioning drive by means of a linear motor, can be arranged along the conveyor path. Furthermore, a transfer apparatus may, however, also be used, for example, for transportation of goods or people, preferably in decentralized arrangements, over long distances and in network systems (transport systems of widely different types).\nIn conventional transfer apparatuses, for example having 3-phase motors, each phase group is hard-wired to a rigidly associated regulator. Motors arranged in series are regulated via an appropriately associated regulator in a respectively defined drive and transmitter zone. However, systems such as these have the disadvantage that every secondary part and every workpiece carrier must be associated with one fixed zone of the primary part, in order that it can be moved deliberately along the primary part. In some circumstances, this means that relatively long distances must be maintained between the individual secondary parts and workpiece carriers. In addition, all the secondary parts are moved along the primary part in accordance with a fixed predetermined regulation rule. In general, this allows only relatively poor positioning accuracy to be achieved.\nIt would therefore be desirable and advantageous to provide an improved transfer apparatus to obviate prior art shortcomings and to attain an increase in throughput while yet improving positioning accuracy."}
{"text": "3GPP Release 12 specifies Proximity-based services (ProSe) (see, for example, Non-patent Literature 1 and 2). The ProSe includes ProSe discovery and ProSe direct communication. The ProSe discovery identifies that radio terminals capable of performing ProSe direct communication (i.e., ProSe-enabled UEs) are in proximity of each other. In an example, the ProSe discovery can be performed through a procedure in which a ProSe-enabled UE detects another ProSe-enabled UE by using only capability of a radio communication technology (e.g., Evolved Universal Terrestrial Radio Access (E-UTRA) technology) possessed by these two UEs. In another example, the ProSe discovery can be performed by a radio access network (E-UTRA Network (E-UTRAN)) or a core network (Evolved Packet Core (EPC)).\nThe ProSe direct communication enables establishment of communication paths between two or more ProSe-enabled UEs that are in direct communication range after the ProSe discovery procedure is performed. In other words, the ProSe direct communication enables a ProSe-enabled UE to communicate with another ProSe-enabled UE directly without communicating through a base station (eNodeB). The ProSe direct communication may be performed by using a radio communication technology that is also used to access a base station (eNodeB) (i.e., E-UTRA technology) or by using a wireless local area network (WLAN) radio technology (i.e., IEEE 802.11 radio technology).\nIn 3GPP Release 12, a ProSe function communicates with a ProSe-enabled UE through a Public Land Mobile Network (PLMN) and assists the ProSe discovery and the ProSe direct communication. The ProSe function is a logical function that is used for PLMN-related operations required for the ProSe. The ProSe function provides functionality including, for example, (a) communication with third-party applications (ProSe Application Server), (b) authentication of UEs for ProSe discovery and ProSe direct communication, (c) transmission of configuration information (e.g., designation of radio resources and transmission power) for ProSe discovery and ProSe direct communication to UEs, and (d) provision of EPC-level ProSe discovery. The ProSe function may be implemented in one or more network nodes or entities. In this specification, one or more network nodes or entities that implement the ProSe function are referred to as a “ProSe function entity” or a “ProSe function server”.\nThe ProSe direct communication in 3GPP Release 12 is one example of the inter-terminal direct communication. Similarly to the ProSe in the 3GPP Release 12, inter-terminal direct communication in a public land mobile network (PLMN) includes the discovery phase and direct communication phase that are assisted by a function or a node (e.g., ProSe function) located in the network. The inter-terminal direct communication is performed between two or more radio terminals in proximity of each other without communicating through any network node (e.g., a base station). The inter-terminal direct communication is also referred to as “device-to-device (D2D) communication” or “peer-to-peer communication”. The ProSe direct communication is an example of the inter-terminal direct communication and is also referred to as “ProSe communication”.\nThe term “public land mobile network” in this specification indicates a wide-area radio infrastructure network, and means a multiple-access type mobile communication system. The multiple-access mobile communication system enables mobile terminals to perform radio communication substantially simultaneously by sharing radio resources including at least one of time resources, frequency resources, and transmission power resources among the mobile terminals. Typical examples of multiple-access technology include Time Division Multiple Access (TDMA), Frequency Division Multiple Access (FDMA), Code Division Multiple Access (CDMA), Orthogonal Frequency Division Multiple Access (OFDMA), and any combination thereof. The public land mobile network includes a radio access network and a core network. Examples of the public land mobile network include a 3GPP Universal Mobile Telecommunications System (UMTS), a 3GPP Evolved Packet System (EPS), a 3GPP2 CDMA2000 system, a Global System for Mobile communications (GSM (Registered Trademark))/General packet radio service (GPRS) system, a WiMAX system, and a mobile WiMAX system. The EPS includes a Long Term Evolution (LTE) system and an LTE-Advanced system.\nPatent Literature 1 discloses that in response to detecting congestion in a network or a failure in a network node (e.g., base station), a server located in the network selects a terminal group to switch to inter-terminal direct communication and transmits, to radio terminals belonging to this terminal group, information about permission for direct communication and configuration information for a communication environment for the direct communication. The information about permission for direct communication indicates, for example, an identifier of radio terminals belonging to the terminal group. The configuration information for the communication environment indicates, for example, a frequency band, transmission power, and an identifier of a receiving terminal, which are used for the inter-terminal direct communication."}
{"text": "1. Field of the Invention\nThe present invention relates to key duplication equipment, and more particularly, to a key positioning fixture for positioning a master key and key blank relative to each other and relative to a key follower and key cutter of a key cutting machine.\n2. Description of the Prior Art\nUse of a conventional key cutting machine having a mechanical key follower and a rotating cutter wheel to duplicate single-sided keys which include a bitted surface on only one side of the key blade produces relatively predictable results. In such single-sided key duplicating operations, the linear root or base surface of a master key and a key blank may be easily aligned with a corresponding linear base surface of a vice-like key clamping device on a relatively repeatable basis may be used to accurately achieve the necessary vertical, lateral (or horizontal) and angular (tilt) alignment between a single-sided master key and a single-sided key blank to yield an operable duplicate key. As a result, the frequency of single-sided key miscuts resulting from the use of conventional key positioning fixtures is acceptably low.\nUsing key positioning fixtures available on conventional key cutting machines to duplicate double-sided keys has historically resulted in unacceptably high miscut rates. Double-sided keys are symmetrical with respect to the key longitudinal axis and are typically used for automotive applications. The original or master copy of a double-sided key which is provided by the original equipment manufacturer to the original automobile purchaser is typically mass-produced on high speed, high cost code-based keel duplication systems which yield keys with high precision mirror image bit notch patterns on the opposing, parallel-oriented sides of the key blade.\nThe low cost, mass-produced key cutting equipment typically used by locksmiths, discount stores, hardware stores and related mass-merchandising outlets while functioning adequately to reproduce single-sided keys, experience severe problems in maintaining the necessary critically accurate alignment between the master key to be duplicated and an uncut double-sided key blank. Alignment errors result from the fact that a key cutting machine operator must use trial and error techniques to properly secure both the double-sided master key and the double-sided key blank in their respective key clamping vices to achieve the necessary vertical, lateral and tilt alignment between the two keys by using the flat, horizontal base surface of the key clamping vices. Because both sides of a double-sided master key blade include an undulating bit notch pattern, the necessary critically accurate alignment between the master key and the key clamping vice cannot be reliably achieved.\nA second source of alignment error arises when the key cutting machine operator removes from the key clamping vices both double-sided master key and the double-sided key blank and invents both keys to allow the second side of the double-sided master key to be duplicated on the key blank. Once again the operator through a trial and error process must establish the critical up/down, left-right and tilt alignment between the double-sided master key and key blank. Because now even the lower \"reference\" surface of the partially cut double-sided, key blank includes an undulating bit notch pattern, the problem of achieving the necessary multiple axis alignment between the master key and the key blank has been rendered even more serious.\nBecause of the number of variables involved in performing the side one key duplication process followed by the side two key duplication process, both of which require that precision alignment be achieved between the master key and key blank, the use of conventional key cutting equipment to duplicate double-sided keys has resulted in unacceptably high miscut rates and correspondingly high customer dissatisfaction."}
{"text": "In a computer network environment, a computer user can retrieve information from a local source, such as a local disk drive, or from a remote source such as a file on a remote server. While the user may locally store the files that the user wishes to interact with, it is common to store other files on a server. For example, where numerous users using different computers may need to sequentially access and read or modify a single file, such as in a collaborative workplace, such a file is typically located on a server that is accessible to the various multiple users.\nBecause a user, or the user's computer, most likely does not have precise and timely information regarding the contents of a particular directory, the user's computer typically requests that the directory be enumerated so that the computer may determine what is and is not stored in the directory of interest. For example, an “open” operation directed to a file in the directory will succeed only if the particular file exists in the specified directory.\nThere are many other situations where an application on a computer may make one or more calls to the server for directory information. For example, some applications trigger open or query-directory calls after each letter of a file name is entered. Thus, when the file name “cabbage” is typed into such an application for opening, seven separate open or query calls are sent to the server. The first is sent when the first “c” is typed in, the second when the first “a” is typed in and so forth with “b,” “b,” “a,” “g,” and “e.”\nExcessive querying of the server for directory and/or file information wastes both communication medium bandwidth and server bandwidth. This in turn either decreases the speed with which networked machines are served or decreases the number of machines serviceable by the server. Accordingly, there is a need for a system that minimizes the number of calls to a server for such information, in order to save the resources of both the server and the communication medium."}
{"text": "This invention pertains to the making of ammunition cartridges for rifles such as those used by hunters, military, and competitive shooters.\nRifle chamber and ammunition cartridge designations are standardized by the Sporting Arms and Ammunition Manufacturer's Institute (SAAMI). ANSI maintains a corresponding standard Z299.4. These standards define, among other things, the physical dimensions of the cartridge and chamber for each cartridge designation by way of a mechanical drawing specifying the dimensions and tolerances for each feature. These dimensions and tolerances dictate how the cartridge will fit into the chamber, take into account changes in the cartridge dimensions during firing, and also account for normal manufacturing variation to ensure that all commercial ammunition will function in all commercial rifles. Allowable variations are small but they can have a significant effect on accuracy because they may alter the way a bullet enters the barrel which affects how it leaves the barrel which affects downrange accuracy. One of the motivations for hand loading ammunition is to take advantage of the ability to adjust the final dimensions of the cartridge to closely match the chamber of a particular rifle and to also decrease the variation from cartridge to cartridge thereby increasing accuracy and consistency.\nAmmunition cartridges are assembled from a case, a primer, powder, and a bullet and are put into several broad classes based on the type of case that is used: rimmed, rimless, and belted being the most common types. Each of these three types of cases use a different physical feature on the case to locate the case inside the chamber, which is commonly called ‘headspacing’ but within the SAAMI specification this is called ‘breeching’. Rimmed and belted cases are breeched by (“headspace off of” is the common terminology) the rim or belt, both features being located at the head of the case (end opposite the bullet). Rimless cases are breeched by (“headspace off of”) the shoulder, the conical transition between the larger cylindrical body of the case, which holds the powder, and the smaller cylindrical neck, which holds the bullet. This is a difference that the prior art has not addressed. All of the different reloading dies, overall-length-gauges, bullet comparators, other tools, and the techniques for using them that are contained in the prior art and commercially available make no differentiation between these different case types, essentially treating all of them as if they were of the rimmed type where all critical dimensions are referenced from the head of the case and alignment is controlled by adjustments to the body and neck of the case. For shoulder breeching cartridges critical dimensions are properly referenced to the case shoulder and it is the case shoulder that should be used to align the bullet.\nAmong other things, careful hand loaders are concerned with the distance that the bullet moves before it engages the rifling of the barrel (called ‘bullet jump’) and the concentricity of the bullet with the bore of the barrel. Several problems arise when loading shoulder breeching cartridges with existing die designs. One is that variations in the length of the case body between the head and the shoulder are translated into variations in bullet jump. This does not occur with rimmed or belted cases because the rifle chamber and the loading press/die are using the same control point, the rim or belt. Existing die and presses are designed to control the Cartridge Overall Length (C.O.L.), the length from the head of the case to the tip of the bullet. This is the sum of two values: the length of the case body from the head of the case to the shoulder and the length from the shoulder to the bullet tip. With existing die designs, a reduction in the length of the case body will result in an increase in the distance from the shoulder to the bullet tip but the C.O.L will remain unchanged. For rimmed and belted cases it is changes in C.O.L. that cause changes in bullet jump so these case designs are not affected by changes in body length. But in a shoulder breeching case, it's changes in the distance from the shoulder to the bullet that cause changes in bullet jump. Therefore, with present die designs any change in the case body length will cause a change in the final shoulder to bullet length, which will cause a change in bullet jump. Because variation in body length from case to case naturally exists, uncontrolled variation in bullet jump must be the result with existing designs.\nAnother problem that arises from using existing die designs with shoulder breeching cartridges is that changes in bullet shape also affect bullet jump. Even within a box of the same bullets from a reputable maker significant amounts of variation are expected, so many loaders weigh each bullet to create batches that all have the same weight. Any change in weight necessarily translates into changes in length at or near the tip, so even when using the same bullets from the same box there is variation in length that will cause changes in bullet jump. Additionally, hand loaders use a variety of bullets in the same case that have significant differences in the shape of the tip. A longer bullet will cause the bullet to be pushed farther into the case, which increases bullet jump. Therefore, a hand loader wishing to control bullet jump needs to carefully adjust existing seating die designs every time they use a different bullet type.\nExisting seating die designs also handle concentricity incorrectly for shoulder breeching cases. A bullet is deformed by the rifling grooves as it is forced into and down the barrel and those grooves force the bullet to spin about the central axis of the barrel (not the central axis of the bullet). If the central axis of the bullet is not perfectly concentric to the axis of the barrel then the bullet will be asymmetrically deformed and the center of mass will be forced to rotate about the axis of rotation. The result is a bullet which wobbles in flight. This has an unpredictable, and therefore deleterious, effect on downrange accuracy.\nThe problem with the existing tools and techniques appears to arise from a lack of appreciation of what the shoulder is doing in a shoulder breeching cartridge case. For rimmed and belted cartridge cases the shoulder is just a transition between the enlarged body and the neck and it does not contact the shoulder of the chamber. Rimmed and belted cases rely on the interface of the rim and belt with the face of the rifle chamber to control the depth of the bullet and on the fitment of the body with the walls of the chamber to control concentricity. Therefore, in rimmed and belted cases the shoulder doesn't contribute to accuracy and can be ignored. Careful hand loaders, particularly competitive shooters using what is commonly called ‘bench rest’ loading techniques and tools (which do not take into account what the shoulder is doing) therefore spend a great deal of time worrying about aligning the exterior of the neck (which positions the bullet) to the body in an effort to control concentricity of the bullet and the bore of the barrel. However, in a shoulder breeching cartridge the shoulder of the case is pushed against the shoulder of the chamber and is the only point of contact between the case and chamber. This forced contact between case and chamber at the shoulder is what aligns the case and controls bullet depth. The body does not contact the chamber walls so it is the body that should be ignored, not the shoulder.\nThis is a critical distinction that is being ignored by the existing technology so I have designed a set of tools for loading cartridges that are similar to existing tools but are designed specifically for shoulder breeching rifle cartridges. These include tools for measuring the chamber and ammunition properly, a case hone, a neck sizing die, and a bullet seating die, each of which is the subject of a separate invention disclosure. This disclosure is for the bullet seating die.\nI will be using the term ‘datum line’ in the description of the invention. This is a term I have found used in some of the relevant literature and it appears to have been in use and well accepted for a long time and thus reasonable to adopt. However, the only official documentation is the SAAMI specification, which does not use this term, so it needs to be defined. The ‘datum line’ is not a visible feature on the cartridge but is the median diameter of the conical shoulder. In the SAAMI Cartridge and Chamber Drawings for shoulder breeching cartridges the datum line is the line drawn through the center of the shoulder, denoted with a “B” for being a ‘basic dimension,’ and also denoted with an ‘x’ inside an ‘o’ on the chamber drawing for being the ‘headspace dimension’. This is also the only place in the drawings where the dimension for the cartridge and the matching dimension for the chamber are the same value. This means that the shoulder is the only place where the cartridge is expected to make contact with the chamber. For an example, see the drawings for the 30-06 Springfield and find the line through the shoulder with a diameter of 0.375 inches.\nI need to create a new term: ‘bullet datum’, which I define as the place where the cross sectional diameter of the bullet is the same as the bore diameter specified in the SAMMI Cartridge and Chamber Drawing. For the 30-06 Springfield this value is 0.300 inches. This is similar to the ‘ogive’, which refers to the forward most part of the bullet having the base diameter; this is 0.309 inches for the 30-06 Springfield. The bullet datum will always be slightly forward of the ogive. Both are near the middle of the bullet and neither are visible features. No one else uses this part of the bullet as a reference point so this is a unique concept. I use it because the exact position of the ogive is difficult to locate, so an alternative is needed. Almost half of a typical bullet has the ogive diameter so it's hard to know exactly where on the bullet a measurement taken from the ogive diameter is being taken from. However, there is only one place on the bullet which has the bore diameter so it's very easy to precisely define and that place is fully on the taper of the bullet so it's pragmatically easier to locate precisely.\nThis disclosure does not discuss how to measure a chamber in order to know what the proper distance is to set up the internal assembly of this die, nor how to measure a cartridge that is produced by this die to ensure that it is properly sized. The technology for making these measurements precisely does not exists. The challenge in both situations is that the reference point is in the middle of a conical section so there isn't a physical feature to anchor one end of a measurement device on nor any visible indicator of where it is. It is still probably possible to make this die work well through trial and error but it would be better to solve the problem on how to take a precise measurement from the center of a cone. I have developed an elegant solution to the problem but because the result is a tool of general applicability and not just a tool for measuring rifle chambers and setting up this seating die it is the subject of another disclosure."}
{"text": "As the manual type spraying container, a foaming container or a container for injecting a liquid as it is, there is commercially sold a liquid injection container which comprises: a container body; a cylinder depending into the container body from a cap-shaped member fitted on the outer face of the mouth/neck portion of the container body; and an actuator having a stem erected through the top wall of the cap-shaped member from a cylindrical piston, which is fitted in the cylinder, and a head fitted on the upper end of the stem and having a nozzle opened forward. The actuator is moved up and down with respect to the cylinder for sucking a liquid in the container body into the cylinder and discharging the liquid in the cylinder out of the nozzle through the stem. Also sold is a container which is enabled to inject one of those spray, foam and liquid selectively by switching a liquid injection passage formed in the head.\nIncidentally, these containers have different nozzle structures but are substantially similar except the structure of the head having the nozzle.\nIn the aforementioned liquid injection containers of the prior art, the actuator composed of the cylindrical piston, the stem and the head is biased upward by the spring so that it is raised back by the spring after it has been depressed. As a result, the depression of the actuator is accomplished against the biasing force of the spring so that it requires a considerable force.\nThe present invention contemplates to provide a liquid injection container, from which the spring is eliminated and in which the finger knob is added to the top face of the head of the actuator, so that the actuator may be moved up and down by the finger inserted into the finger knob. These actuations are facilitated by erecting auxiliary walls along the two sides of the head of the actuator for acting as a finger grip and guiding the head."}
{"text": "The present invention relates to wet shaving razors and, more particularly, to a disposable wet shaving razor having a \"rubberized\" handle for a more comfortable and secure grip.\nDisposable wet shaving razors which are intended to be discarded after the shaving blade becomes too dull to provide an acceptable shave are in widespread use. In order for a disposable razor to be commercially feasible, its manufacturing cost must be kept at a minimum.\nThe desirability of providing a wet shaving razor with a rubberized handle which provides a more comfortable and secure grip, especially when held with wet and sometimes soapy hands, has been recognized. However, known techniques for fabricating wet shaving razors with rubberized handles have the drawback of being too costly to use in the manufacture of disposable razors. For example, U.S. Pat. No. 4,949,457 to C. J. Burout III discloses and claims a composite razor handle having a rigid inner core of thermoplastic material and a molded flexible, resilient covering layer. The flexible covering layer is formed by a separate molding process and attached to the rigid inner core, which is provided with elongated recesses to prevent the rotational slippage of the covering layer and notches for anchoring the covering layer to prevent it from peeling away from the core. Although the Burout patent states that the composite razor handle taught therein may be used in disposable as well as nondisposable razors, the additional molding and attachment steps required to form the flexible, resilient covering layer for the Burout razor handle would add significantly to the manufacturing cost of a razor. Therefore, use of the Burout razor handle in a disposable razor is disadvantageous, if not impractical.\nAccordingly, a need clearly exists for a disposable razor having a rubberized handle which can be manufactured at low cost."}
{"text": "Plaque is a common factor in caries, gum disease and discoloration of teeth and greatly contributes to their development. Plaque is initiated when cariogenic bacteria adhere to pellicle, a proteinaceous film on the surface of teeth. Plaque, in turn, acts as a nucleus for the formation of calculus. As calculus matures and hardens it tends to stain due to the absorption of dietary chromagens. Additionally, oral restorative materials may be inherently susceptible to build-up of stain from dietary chromagens. It is desirable to have a means to avoid stain absorption and adherence of bacteria to hard tissue and surfaces of the oral environment.\nSilicone oils, because of their hydrophobic nature, have been suggested for inclusion in dentrifices to inhibit the staining process. However their adhesion and retention on tooth surfaces is typically quite low.\nU.S. Pat. No. 5,078,988 to Lin et al. discloses dentifrices including modified aminoalkyl silicones. The modified silicones are said to form a hydrophobic layer on the teeth for prevention of caries and stain. PCT patent application number WO 91/13608 to Rolla et al. discloses dentifrices comprising a liquid silicone oil and a fat-soluble antibacterial agent, which is described as being useful for protection of teeth against plaque formation due to a slow release of antibacterial agent into the saliva.\nU.S. Pat. No. 4,981,903 to Garbe et al. discloses pressure-sensitive or non-pressure sensitive adhesive compositions comprising a vinyl polymeric backbone with grafted pendant siloxane polymeric moieties. These compositions are disclosed to be useful as good topical application binding materials for application in cosmetics and medicaments, and also as sealant compositions for porous materials such as paper and wood. See col. 3, lines 25 through 31. U.S. Pat. No. 4,972,037 to Garbe et al. discloses compositions useful as an adhesive at room temperature which comprise a copolymer having both pendant fluorochemical groups and pendant polysiloxane grafts. These compositions are also useful in topical applications, including the application of cosmetics and medicaments and for sealant compositions for porous materials such as paper and wood. See col. 3, lines 25 through 33. U.S. Pat. No. 4,981,902 to Mitra, et. al. discloses non-pressure-sensitive adhesive acrylate or methacryate polymers having pendant polysiloxane grafts. The polymers comprise monomers having polar functionality, and are described as useful in coating compositions for animal bodies.\nU.S. Pat. No. 4,693,935 to Mazurek discloses pressure-sensitive adhesive compositions comprising a copolymer having a vinyl polymeric backbone having grafted thereto polysiloxane moieties. U.S. Pat. No. 4,728,571 to Clemens et al. discloses release coating compositions comprising polysiloxane grafted copolymer and blends thereof on sheet materials."}
{"text": "1. Field of the Invention\nThe invention applies to wrapping for used chewing gum and provides easy hygienic and ecological disposal of used chewing gum which improves the environment.\n2. Description of Prior Art\nUsed chewing gum, thrown out in public places, for example on pavements, floors of public buildings, or on grass plots, contributes to the pollution of the environment and substantially diminishes the utility and aesthetic value of public places. Removal of used chewing gum is very problematic, because the gum usually firmly sticks to surfaces upon which they lie.\nOne known approach for removal of sticky chewing gum from pavements or floors is to use various machines. Typical machines use a dissolving agent or use local cooling to make mechanical releasing easier. High price and high operating costs are among the disadvantages of using such machines. Yet another disadvantage is that cleaning is performed only from time to time.\nTo try to prevent the above mentioned problems, keepers or owners of publicly open places often use hygienic devices such as small containers, in which an exchangeable insert is pushed in. One disadvantage of those devices is that their usage is limited to placement within buildings, such as banks, hotels, restaurants and the like. This approach does not solve the problem of discarding of used chewing gum in open areas, on pavements and grass plots.\nAnother known solution is a device that includes many pieces of wrapping paper. Each paper can be ripped off and used as a wrap for a piece of used chewing gum. Disadvantages of this solution are complicated handling, the necessity of using both hands to open the wrapping paper, the need to use bare hands or putting paper on to consumer's lips, paper that may have already been contaminated by having been carried for a long time in one's pocket or in one's bag."}
{"text": "The present disclosure relates to an optical apparatus and its vibration removing method.\nFor a lithography technology that plays an important role in the miniaturization of semiconductor devices, ArF lithography in which an ArF excimer laser having an exposure wavelength of 193 nm is used as an exposure light source is currently used for mass-production. Further, an immersion technique (also called “ArF immersion lithography”) for increasing the resolution of an exposure device by filling a space between its objective lens and a wafer with water has also begun to be used for mass production. In order to advance the miniaturization of semiconductor devices even further, development of various techniques for putting EUVL (Extremely Ultraviolet Lithography) having an exposure wavelength of 13.5 nm to practical use has been in progress.\nIn an apparatus using EUV light, an optical system is disposed inside a vacuum chamber so that an optical path is in a vacuum. When there is a vacuum pump and/or a driving mechanism such as a stage, vibrations occur. Further, in some cases, vibrations occur in a light source. Therefore, an exposure apparatus using a vibration removal table is disclosed.\nAn EUV exposure apparatus disclosed in Japanese Unexamined Patent Application Publication No. 2005-136120 includes two vibration removal tables. While one of the vibration removal tables removes vibrations from a stage apparatus and a loader, the other vibration removal table removes vibrations from an illumination optical system and a projection optical system. A vibration removal apparatus disclosed in Japanese Unexamined Patent Application Publication No. H11-264444 includes a vibration removal table using an air spring.\nAn exposure apparatus disclosed in Japanese Unexamined Patent Application Publication No. H11-327155 includes an exposure apparatus main unit, a laser oscillation unit, and a vibration removal table. All the units disposed in the exposure apparatus main unit are mounted on the vibration removal table. Further, the laser oscillation unit is connected to the exposure apparatus main unit through an accordion connection part. In this way, vibrations in the laser oscillation unit are prevented from propagating to the exposure apparatus main unit."}
{"text": "Software running on a processor, microprocessor, and/or processing unit to provide certain functionality often changes over time. The changes can result from the need to correct bugs, or errors, in the software files, adapt to evolving technologies, or add new features, to name a few. In particular, embedded software components hosted on mobile processing devices, for example mobile wireless devices, often include numerous software bugs that require correction. Software includes one or more files in the form of human-readable American Standard Code for Information Interchange (ASCII) plain text files or binary code. Software files can be divided into smaller units that are often referred to as modules or components.\nPortable processor-based devices like mobile processing devices typically include a real-time operating system (RTOS) in which all software components of the device are linked as a single large file. Further, no file system support is typically provided in these mobile wireless devices. In addition, the single large file needs to be preloaded, or embedded, into the device using a slow communication link like a radio, infrared, or serial link.\nObstacles to updating the large files of mobile processing devices via slow communication links include the time, bandwidth, and cost associated with delivering the updated file to the device. One existing solution to the problem of delivering large files to mobile processing devices includes the use of compression. While a number of existing compression algorithms are commonly used, often, however, even the compressed file is too large for download to a device via a slow, costly, narrowband communication link.\nAnother typical solution for updating files uses difference programs to generate a description of how a revised file differs from an original file. There are available difference programs that produce such difference data. However, as with compression, the difference files produced using these difference programs can sometimes be too large for transfer via the associated communication protocols."}
{"text": "For Windows™ or like platform users, developing mainframe applications usually entails manually uploading a program from the Windows™ or like platform, logging on to the mainframe or host computer systems, and submitting the job, that is, executing the program, and waiting for the output. For example, CA-Easytrieve, which includes fourth generation language (4 GL) query and reporting capabilities and allows users to access many IBM mainframe data formats, requires the user to manually perform those tasks."}
{"text": "1. Field of the Invention\nThe present invention relates to litter handling devices and more particularly pertains to a litter filtering system for filtering animal litter to remove impurities therefrom.\n2. Description of the Prior Art\nThe use of litter handling devices is known in the prior art. More specifically, litter handling devices heretofore devised and utilized are known to consist basically of familiar, expected and obvious structural configurations, notwithstanding the myriad of designs encompassed by the crowded prior art which have been developed for the fulfillment of countless objectives and requirements.\nKnown prior art litter handling devices include U.S. Pat. Nos. 5,211,133; 5,272,999; 5,220,886; 5,178,099; and 5,032,254.\nWhile these devices fulfill their respective, particular objectives and requirements, the aforementioned patents do not disclose a litter filtering system for filtering animal litter to remove impurities therefrom which includes first and second containers operable for receiving a supply of litter for use by a pet, and a screen assembly adapted for positioning over the first container and receiving the second container in an inverted position on top of the first container such that the screen assembly is positioned between the containers, wherein the assembly of the containers and screen assembly can be manually rotated so as to position the first container in an inverted position, whereby the litter is caused to filter through the screen assembly and into the second container.\nIn these respects, the litter filtering system according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in so doing provides an apparatus primarily developed for the purpose of filtering animal litter to remove impurities therefrom."}
{"text": "Semiconductor devices are used in many electronic and other applications. Semiconductor devices comprise integrated circuits that are formed on semiconductor wafers by depositing many types of thin films of material over the semiconductor wafers, and patterning the thin films of material to form the integrated circuits.\nThere is a trend in the semiconductor industry towards reducing the size of features and/or improving performance of the semiconductor devices. For example, features of devices are scaled to improve on current performance, performance of parasitic resistances are decreased etc. However, such technological progress requires overcoming many challenges. One challenge involves the contact formation in the front-end-of-line that includes silicidation of active area and forming contacts to it through an insulating layer. Scaling challenges these processes by decreasing the contact sizes as well as the contact to contact spacing. Increasingly silicidation introduces defects that significantly reduce process yields. Another challenge is the improvement in the performance of diodes used in semiconductor devices.\nSchottky diodes, which are metal-semiconductor diodes, are conventionally formed by contacting silicide regions with semiconductor regions. This is because of the compatibility of silicide based diodes with conventional semiconductor processing. However, such diodes have a number of limitations. For example, they are susceptible to increased leakage currents at corners, for example, due to thinner silicidation as well as due to the existence of increased field regions.\nFIG. 1, which includes FIGS. 1a and 1b, illustrates a conventional silicide based Schottky diode, wherein FIG. 1a illustrates a top view and FIG. 1b illustrates a cross-sectional view of the semiconductor device.\nReferring to FIG. 1b, a Schottky contact is made between a silicide contact region 23 and a first doped region 20. The silicide contact region 23 is formed on the first doped region 20 and is formed between adjacent isolation regions 50. The contact to the semiconductor is made through silicide regions 22 formed over the second doped region 21, which is a heavily doped region for contacting the semiconductor portion of the diode.\nThe isolation regions 50 prevent the silicide contact region 23 from directly contacting the silicide regions 22. In some devices, the silicide contact region 23 is formed under the contact plugs 31, for example, in trenches under the contact plugs 31, so that a plurality of silicide contact regions 23 are formed under the contact plugs 31."}
{"text": "The rise of the Internet and networking technologies has resulted in the widespread transfer of code, data and files between computers. This material is not always what it seems to be. For example, code that is accessed on a remote machine and downloaded to a computer system can contain hostile algorithms that can potentially destroy code, crash the system, corrupt code or worse. Some of these hostile algorithms are viruses, worms, and Trojan horses.\nHostile, malicious and/or proscribed code, data and files (“code” as used hereinafter generally includes “text,” “data” and “files”) can infect a single computer system or entire network and so posit a security risk to the computer system or network. The user and/or administrator (generally referred to hereinafter as “user”) may wish to intercept, examine and/or control such code. The user might also wish to intercept, examine and/or control other code and/or text as well, for example, text which the user does not know to be hostile, but wishes to intercept nonetheless, such as potentially sexually or racially harassing messages, unsolicited messages, etc. This latter type of code is known hereinafter as “predetermined code.”\nOne method of transferring code, instant messaging, has become increasingly popular. Instant messaging may be used for real time text exchanges, as well as other exchanges, such as image transfer, voice and/or video chat, interactive games, code transfers, remote assistance, whiteboarding, code and/or application sharing, etc. Thus a user may transmit code through instant messaging as well as text messages.\nThe transmission of code through instant messaging may include hostile, malicious, predetermined and/or proscribed code (generally referred to hereinafter as “proscribed code.”) Additionally, although messages themselves are usually clear text, encryption may be used on messages and/or code etc., e.g. through Secure Sockets Layer (SSL) and Secure Multi-Purpose Internet Mail Extension (S/MIME).\nAs with other transfers, it may be desired to secure instant messaging. However, instant messaging applications, such as AOL Instant Messenger (AIM), NET Messenger (including Windows Messenger, MSN Messenger, etc.), Yahoo Messenger, etc., may be difficult to secure for a number of reasons. For example, in an enterprise installation, IM applications may be installed illegitimately, and thus serve as a conduit for proscribed code. Even those installations that are legitimate may provide, through IM, proscribed code to the user's machine.\nIndeed, since Instant Messaging protocols may make use of tunneling and port scanning in order to transmit messages, they may transmit proscribed code without detection. Thus, FTP, Telnet and/or HTTP ports for example may be used by an IM application. As these ports are often left open for other applications, it is difficult to shut them off for IM applications.\nSimilar problems exist with peer-to-peer and peer-to-peer like applications, e.g., Gnutella (peer-to-peer), Kazaa (peer-to-peer like, as it utilizes a central server.) (Both are referred to herein as “peer-to-peer” or “P2P.”) These applications may be installed illegally or without authorization on a system, and may make use of tunneling, port scanning and other techniques making detection difficult.\nTherefore, it would be beneficial to have apparatus, methods and articles of manufacture to simply and effectively intercept, control, and/or examine incoming and outgoing instant messaging and/or peer-to-peer code in an efficient and effective manner transparently or almost transparently to the end-user, with little or no operational effort required by the user.\nIt would further be beneficial to have apparatus, methods and articles of manufacture to simply and effectively intercept, control, and/or examine incoming and outgoing secure instant messaging and/or peer-to-peer code in an efficient and effective manner transparently or almost transparently to the end-user, with little or no operational effort required by the user."}
{"text": "With the increasing popularity of mobile devices, device manufacturers have sought ways to create devices with more processing power and other capabilities in a smaller size. In many cases, increasing the processing power of a device through the use of more powerful processors or adding other components may lead to an increase in heat generated by those components inside the device. Also the decreasing size of a device may inhibit the device's ability to dissipate heat. Accordingly, smaller and more powerful devices may be more prone to localized extreme temperatures. For example, a hot spot may develop which is uncomfortable to a user or which may result in damage to the component.\nCertain implementations and embodiments will now be described more fully below with reference to the accompanying figures, in which various aspects are shown. However, various aspects may be implemented in many different forms and should not be construed as limited to the implementations set forth herein. Like numbers refer to like elements throughout."}
{"text": "This application is related to Japanese Patent Application No. H11-295748, filed Oct. 18, 1999, the disclosure of which is incorporated by reference herein in its entirety.\n1. Field of Invention\nThe present invention relates to head protective bags for vehicle passengers, and, more particularly, to bags inflated along a side window or the like when a side of the vehicle is crashed or the vehicle is rolled over. More specifically, the present invention relates to head protective bags for vehicle passengers that are inflated by gas. Further, the present invention relates to head protective devices for vehicle passengers having the head protective bags and vehicles having the protective devices.\n2. Description of the Related Art\nThis type of head protective bag for vehicle passengers is described in WO 96/26087 and is illustrated specifically in FIGS. 1 and 9 of that publication. This known bag includes a duct portion (main airway) extending along an upper edge of the bag from an end (a front end) of the bag near the vehicle\"\"s front end to an end (a rear end) of the bag near the vehicle\"\"s rear end and a number of cell portions (small cells) connected to the duct portion are extending downward. The head protective bag described in WO 96/26087 will now be described with reference to FIG. 6 of this application.\nThis head protective air bag device 100 has main components including an elongated duct 106, an air bag body 112, an inflator 116, and a belt-like strap 118. The duct 106 extends along a front pillar portion 102 and a roof side rail portion 104. The air bag body 112 is accommodated in the duct 106 in a folded state and is secured to a vehicle body at a front securing point 108 and a rear securing point 110. The inflator 116 is connected to a rear end of the duct 106 through a hose 114 and generates gas when a predetermined load acts on a side of the vehicle body. One end of the strap 118 is secured to the vehicle body, and the other end of the strap 118 is secured to a rear end of the air bag body 112. The air bag body 112 is formed by connecting a plurality of cells 120. Each cell 120 is substantially cylindrical and is located such that a longitudinal direction of the cell 120 corresponds to a substantial upward or downward direction of the vehicle.\nIn this structure, the inflator 116 generates gas when a predetermined load acts on a side of the vehicle. The generated gas is sent to each cell 120 of the folded, air bag body 112 through the hose 114 and the duct 106. Accordingly, the cell 120 is inflated in a substantially cylindrical shape such that its longitudinal direction corresponds to the vehicle s upward or downward direction. The air bag body 112 is thus inflated along a window 122 in a curtain-like shape. Further, since the rear end of the air bag body 112 is connected to the vehicle body through the strap 118, the rear end of the air bag body 112 is reliably located in an upper section of a center pillar portion 124.\nIt is an objective of the present invention to provide a head protective bag for vehicle passengers, a lower section of which is tightly expanded when the bag is inflated, and a protective device and a vehicle having this protective bag.\nA head protective bag for vehicle passengers of the present invention is located near a corner defined by a ceiling and a side wall of a passenger compartment and is inflated downward along the side wall by means of gas introduced from a gas inlet. The bag includes a front coupling portion secured to a first pillar near a front side of a vehicle and a rear coupling portion secured to a second pillar located rearward from the first pillar. The bag is characterized in that a dimension L between a lowermost front coupling portion and a lowermost rear coupling portion when the bag is maximally inflated (without being installed in the vehicle) is smaller than a distance S from a position at which the front coupling portion is secured to the first pillar to a position at which the rear coupling portion is secured to the second pillar.\nWhen the head protective bag for vehicle passengers is maximally inflated without being installed in the vehicle, the dimension L between the front coupling portion and the rear coupling portion is smaller than the distance S. Accordingly, when the bag is inflated as installed in the vehicle, an increased tensile force is produced between the front coupling portion and the rear coupling portion of the bag. The portion of the bag between these coupling portions is thus tightly expanded. As a result, the tightly inflated bag receives a passenger who hits the bag.\nTo reduce the dimension L with respect to the distance S, it is preferred that a plurality of small cells are arranged in a forward or rearward direction in the bag and that a width of a lower section of the bag in the forward or rearward direction is larger than a width of an upper section of the bag in the forward or rearward direction. Further, to reduce the dimension L with respect to the distance S, the bag may include a cabin side sheet and a window side sheet. A cell in which the gas is introduced is formed between the sheets. A coupling portion is provided in the cell for coupling the sheets together. The number of the coupling portions in an upper section of the cell is larger than the number of the coupling portions in a lower section of the cell.\nIn the present invention, a single inflatable cell may be formed between the sheets."}
{"text": "A pump and, in particular, the armature of the stator of the electrical drive motor of such pumps are usually designed for providing a certain head H0 at zero flow rate with maximum efficiency. Normally, the maximum efficiency is reached by a design that allows providing the required head H0 at zero flow rate in a non-filed weakening mode at a maximum modulation index close to 1.\nEP 2 133 991 describes a speed controlled pump which makes use of field weakening to optimise electrical power consumption at low flow rates."}
{"text": "This invention relates to a roller clamp used to control liquid flow through a flexible tube. The invention finds particular application in the medical field for controlling flow, for example, through I.V. tubing.\nIt is important accurately to control the flow of I.V. fluid to a patient and this is commonly achieved by means of an adjustable clamp on tubing which delivers the fluid to the patient from an I.V. bag or the like. In use, the clamp is opened or closed on the tube to control the rate of flow.\nThe prior art is replete with different designs of tube clamps for the above purpose. One common form of clamp comprises an elongate channel-like housing in which is mounted a wheel or roller on an inclined track which extends lengthwise of the channel. The tubing is clamped between the wheel and the base of the channel, and the flow rate through the tubing is varied by moving the wheel along the inclined track so as to increase or decrease its distance from the base of the channel thereby adjusting the degree to which the wheel exerts pressure on the tubing.\nTubing clamps of the above kind may suffer certain disadvantages namely, because of the linear motion of the wheel along the channel to adjust the flow rate, there may not be a positive positioning of the wheel and should the clamp inadvertently be moved along the tube, the flow rate may be altered.\nThe present invention, therefore, seeks to provide an improved tubing clamp generally of the above type."}
{"text": "The present invention relates to an air bleed system for a compressor of a gas-turbine engine.\nMore specifically, the present invention relates to a centripetal airbleed system positioned between oppositely opposed first and second disks of a gas-turbine engine compressor. The air bleed system includes an annular support affixed to the first and second disks, and a plurality of airbleed tubes mounted in a substantially radial manner within orifices defined along the periphery of the support. Each tube is arranged with a vibration reduction device that reduces vibrations during operation of the engine.\nAn air bleed system is proposed in U.S. Pat. No. 5,472,313 wherein a damping tube is inserted within the radially inner part of each airbleed tube. At a radially outer portion, the damping tube is arranged with longitudinal slots defining axial strips which, during gas turbine engine operation, are centripetally urged against the inside wall of the bleed tube. Friction between the inner damping tube strips and the outer airbleed tube dissipates deformation energy and therefore lowers the dynamic stresses in the outer tube.\nThrough-holes are present at the end of the longitudinal slots and constitute stress concentration zones. This feature requires additional and intricate machining. Further, the service life of such inner damping tubes is less than 100,000 cycles.\nThe mode 1F frequency of the centripetal airbleed-tube of the GE turbojet engine 90-115B is 950 Hz. The aformentioned engine also includes 12 bleed tubes, and the 8F mode frequency operating at 7,125 rpm is (7,125xc3x978/60=) 950 Hz. This is the same frequency as for the 1F mode.\nOne object of the present invention is to lower the dynamic stress in the airbleed tube by dissipating the energy of deformation.\nAnother object of the present invention is to propose an airbleed system based on centrifugal outflow wherein the device damping the airbleed tube vibrations allows for substantially increasing the frequency of the first bending mode, the so-called 1F mode, in order to eliminate crossing of the 1F mode and drive source.\nAccording to the present invention, there is provided a device for reducing the vibrations of an airbleed tube that includes a damping tube maintained within an annular support orifice and which encloses a radially outer part of an airbleed tube. It follows that the radially inner part of the damping tube encloses a portion of the radially outer part of the airbleed tube. An annular gap is defined between the median zone of the damping tube and the airbleed tube.\nThe damping of the centripetal airbleed system of the present invention relies on two principles: flexibility exhibited at the contact between the outer damping tube so that the airbleed tube can act as a spring; and the rigidity at the contact dissipates energy as if it were a spring. Moreover, the outer damping tube may be stressed by bending such that the bending strength of the damping tube enables energy dissipation.\nThe proposed system of the invention allows lowering the dynamic stress in the airbleed tube by dissipating energy of deformation and by an increase in the frequency of the first bending mode, the so-called 1F mode, on account of the increase in equivalent rigidity of the airbleed system.\nIn addition, the proposed outer damping tube is free of slots and therefor holes at the ends of the slots. As a result, the manufacture of the present invention is easier than the machining of a strip-fitted inner damping tube of the prior art.\nIn a first embodiment of the present invention, the radially inner end of the damping tube is constricted sufficiently to contact a peripheral wall of the airbleed tube along a plurality of axial zones defined between the damping and airbleed tubes.\nIn a second embodiment of the present invention, the airbleed tube includes a boss extending from a radially outer surface thereof that engages with the radially inner end of the damping tube. The boss includes a plurality of flat surfaces defined along thereof.\nOther advantages and features of the invention are elucidated in the following illustrative description and in relation to the attached drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a display device. Further, the invention relates to an electronic appliance having the display device for a display portion.\n2. Description of the Related Art\nIn recent years, a thin display device having pixels formed using self-luminous light emitting elements has attracted attention. As a light emitting element, an organic light emitting diode (OLED) or an EL (electroluminescent) element has attracted attention, and have been used for an organic EL display or the like.\nAs a driving method for expressing a multi-gray scale image of a display device using the aforementioned light emitting element, there are an analog driving method (analog gray scale method) and a digital driving method (digital gray scale method).\nThe analog driving method is a method in which current magnitude flowing in a light emitting element is continuously controlled to obtain a gray scale. Whereas, the digital driving method is a method in which a light emitting element is driven by only two states of an ON state (a lighting state with the luminance of approximately 100%) and an OFF state (a state where the luminance is approximately 0%, that is, a non-lighting state).\nNext, brief description is made of an example of a pixel structure of a display device employing the time gray scale method and drive thereof. A circuit shown in FIG. 1 includes transistors 201 and 202, and a light emitting element 203. A gate electrode, a first electrode, and a second electrode of the transistor 201 are connected to a gate signal line 205, a source signal line 204, and a gate electrode of the transistor 202 respectively. A first electrode and a second electrode of the transistor 202 are connected to a power source line 206 and a first electrode of the light emitting element 203 respectively. A second electrode of the light emitting element 203 is connected to a counter electrode.\nNote that it is difficult to define a source electrode and a drain electrode of a thin film transistor (hereinafter referred to as TFT) due to a structure thereof. Here, one of a source electrode and a drain electrode is referred to as a first electrode, and the other is referred to as a second electrode. In general, a lower potential side electrode is a source electrode and a higher potential side electrode is a drain electrode in an n-channel transistor, whereas a higher potential side electrode is a source electrode and a lower potential side electrode is a drain electrode in a p-channel transistor. Accordingly, in the case where there is description concerning a gate-source voltage or the like in description of circuit operation, the aforementioned basis is referred.\nSubsequently, description of FIG. 1 is made with reference to a timing chart in FIG. 2. The source signal line 204 to be selected is determined by an SWE211 (source writing/erasing select signal). Further, the gate signal line 205 to be selected is determined by a GIWE212 (gate writing select signal) and a G2WE213 (gate erasing select signal). Whether the light emitting element 203 emits light or no light is determined by signals of the source signal line 204 and the gate signal line 205. Here, as for an arbitrary wiring, a digital signal “1” is referred to as H (High level), whereas “0” is referred to as L (Low level). It is to be noted that “0” means not only a ground potential but a common potential. A state where a potential is higher than an arbitrary threshold voltage may be H, whereas a state where a potential is lower than an arbitrary threshold voltage may be L.\nBlack is written when the source signal 214 is H. However, if the gate signal 215 is not H at that time, such data is not reflected to the light emitting element 203. Meanwhile, white, that is, data is written when the source signal 214 is L. However, if the gate signal 215 is not H, such data is not reflected to the light emitting element 203.\nSubsequently, the digital driving method is described. With the digital driving method alone, only 2 gray scales can be expressed. Therefore, it is suggested that the digital driving method be used in combination with a driving method for expressing multi gray scales, such as an area gray scale method or a time gray scale method. The area gray scale method is a method in which gray scale is expressed depending on the size of a light emitting area of a sub-pixel provided in a pixel (for example, see Patent Document 1). Further, the time gray scale method is a method in which gray scale is expressed by controlling a light-emitting period and light-emitting frequency (for example, see Patent Documents 2 and 3).    [Patent Document 1] Japanese Published Patent Application No. H11-73158    [Patent Document 2] Japanese Published Patent Application No. 2001-5426    [Patent Document 3] Japanese Published Patent Application No. 2001-343933"}
{"text": "1. Field of the Invention\nThe present invention relates to building panels, particularly for the construction of houses. The invention is particularly concerned with a framework or skeleton panel which is formed as a grid of intersecting members, usually of wood-based material. The term \"panel\" will be used herein to include such a panel framework, and does not imply that this is a finished panel with insulation and/or facing sheets.\n2. Prior Art\nIt is known to construct houses, and other buildings, from panels which are factory made, and which usually contain insulation. Examples of patents showing such panels are as follows:\nU.S. Pat. No. 4,671,032, which issued Jun. 9, 1987 to Reynolds;\nU.S. Pat. No. 4,765,105, which issued Aug. 23, 1988 to Tissington et al.;\nU.S. Pat. No. 4,894,974, which issued Jan. 23, 1990 to Mayhew et al.\nU.S. Pat. No. 5,157,892, issued Jun. 9, 1987 to Ryther; and\nU.S. Pat. No. 5,167,700, issued Apr. 8, 1997 to Wright et al.\nApplicant has had considerable experience both with the panels of the last-mentioned '700 patent, of which he is co-inventor, and with those of the Tisssington et al. patent. Both these patents are concerned with factory made panels having a framework made of wood members, and insulated with rigid foam insulation injected between the wood members in the factory. In the last-mentioned patent, dimensional lumber is used to provide strength, while in the Tissington et al patent it is preferred to use board such as oriented strand board (OSB). OSB is cheaper than standard dimensional lumber and is more resistant to warping. These panels have been used to construct over one thousand buildings of many different types, with great success.\nPanels of the type shown in the Tisssington et al. patent and in the '700 patent have great potential in export markets, since buildings can be erected with these panels using largely unskilled labor. However, it has become apparent that a serious drawback of these and similar panels, i.e panels fully assembled and insulated in the factory, is that the bulkiness of the insulation leads to high shipping costs which are disadvantageous for export markets. The shipping costs could be much reduced if it were possible to ship a kit of structural members which could easily be assembled into a panel on site, without the need to ship the insulation.\nThe present invention accordingly is concerned with a panel in the form of a framework which can be shipped in disassembled form, without any insulation. The basic structural members of the panel occupy about one-fifth of the volume occupied by a fully assembled and insulated, factory produced panel. The structural members of the panel can easily be assembled on site by unskilled workers to form a rigid panel framework. Insulation, and facing sheets for example of plywood, OSB, and many other materials, can be added in accordance with local requirements and availability.\nThe present invention makes use of novel structural members which can easily be put together to form a grid. Panels formed as grids of crossing, interlocking members, are not new per se, and the Mayhew and Ryther patents, as well as the '700 patent aforesaid, show such panels. However, the prior art panels have various drawbacks.\nOne drawback of many of the prior art designs is that the grid forming members, as for example in Mayhew et al. and Ryther, are too thin for their edges to reliably receive nails for the facing sheets.\nAnother drawback has been the need to use structural members additional to those of the basic panel itself. Firstly, such panels have generally needed some kind of rails or studs or side pieces to improve the appearance of any exposed edges of the panels, which otherwise show the ends of the lateral members. Such rails or studs are shown for example at 240 in FIG. 2 of Mayhew, and as 28 in FIG. 1 of Ryther. Secondly, the grid forming members in many cases have not been stiff enough to be used exclusively to form the panel; in Mayhew et al. the members 240 are in the nature of studs which add strength and stiffness. Accordingly, these prior constructions all require the use of members other than the basic grid forming members to complete the panel framework.\nOne special feature of this invention is that a panel framework can be entirely formed from a plurality of structural members of novel design, which can be all identical or can be of only two different designs. These basic members are used not only for the horizontal and vertical members of the panel grid, but also form its four edges, i.e. the side edges and end (e.g. top and bottom) edges of the panel, where they provide smooth relatively uninterrupted edge surfaces, while producing a panel of good strength and rigidity. The panel also has adequate nail receiving surfaces on both sides, even when formed of relatively thin material such as 3/4 inch thick OSB. Insulation and facing sheets can be added after assembly.\nFurthermore, the novel structural members can also be used to produce beams and joists needed for building a house, and also floor panels which need higher strength than the basic wall panels. In fact it is possible to construct houses almost entirely from panels, joists and other parts produced from these one or two types of structural member, along with facing sheets, and insulation if required, which can be obtained locally. This avoids the need to organise shipments of the many different structural members usually needed to make a house, and avoids problems which frequently occur if there is a shortage or breakage of one or two structural members of a specialized design."}
{"text": "The present invention relates in general to electronic memories and in particular to content addressable memory cells and systems and devices using the same.\nThe most prevalent type of memory in data processing and telecommunications applications is the random access memory (RAM). In a RAM, the memory cells are organized in an array of addressable locations, where each location is comprised of one or more cells each storing a bit of data. During a random access, data are associated with an address and the corresponding location in the array accessed (read from or written to). During burst and page accesses, multiple words of data are stored or retrieved from a commonly addressable set of locations in the array, such as those along a single row when the array is organized in rows and columns.\nRandom access memories are not ideal in a number of situations. For example, to search for data stored within RAM, a sequence of addresses must be generated and data from the accessed locations sequentially examined until the search is complete. Typically, there is no guarantee that the data will be found within a predetermined number of access cycles, or even that the data will be found at all. This procedure is inefficient in applications such as networking, ATM (asynchronous transfer mode) switching, voice and image recognition, where time constraints on searching are critical.\nAs a result of the deficiencies of conventional RAM devices, the content addressable memory (CAM) was developed. The CAM is a data associative memory in which data is broadcast to all locations in the cell array at once. A comparison is made in each location between a received piece of data and the data stored in the cells of that location. When a match occurs, a flag is set to an active state based on the associated xe2x80x9cmatchlinexe2x80x9d. The output at the match lines can then be used as data for further processing. For example, the match line output may represent an address in RAM associated with a search target.\nContent addressable memories can be either static or dynamic. The static design is very similar to a static RAM (SRAM) cell, essentially acting as a latch to store each bit of data. The Dynamic CAM, similar to a dynamic RAM, stores data as charge on a capacitor, although charge leakage from the capacitor will cause data deterioration and even loss if left unchecked. Consequently, with dynamic CAMs, the charge on the storage data in the cell capacitors must be periodically refreshed. In conventional dynamic CAMs, time must be taken for cell refresh, generally during which no data accesses can be made to the array.\nAnother problem with conventional dynamic CAMs arises from the fact that during a data access or refresh to the array, match operations cannot normally be performed. This is particularly true with regards to those columns of locations to which data are being read or written during update of the search database. As a result, CAM performance cannot be optimized.\nWhat are needed therefore are CAM cell structures and device architectures which allow simultaneous updates of the stored database while match operations are being performed. Additionally, in the case of dynamic cells, such cells and architectures should retain the reduced power consumption advantage reduced die area a dynamic device provides but appear as a static CAM to an external processing device.\nAccording to one embodiment of the principles of the present invention, a content addressable memory cell is disclosed which includes a first and second storage elements along with first and second transistors for selectively transferring charge between corresponding first and second bitlines and the first and second storage elements. The content addressable memory cell further includes first and second comparelines, each for carrying a corresponding bit of a comparand. A comparator 402 compares first and second bits of the comparand presented on the comparelines with information stored on the first and second storage elements and selectively controls a voltage on a corresponding one of a plurality of matchlines in response.\nAccording to another embodiment of the principles of the present invention, a content addressable memory cell is disclosed which includes first and second pairs of bitlines, along with first and second memory cells. Each memory cell includes a first transistor for coupling the data storage node of that cell to either one of the first pair of bitlines or one of the second pair of bitlines as selected by a signal presented on a corresponding one of first and second wordlines. The memory cell is further associated with a plurality of matchlines for carrying the results of multiple comparisons made against the information stored in the memory cells. Comparator circuitry compares data stored at the first storage node with data presented on a first compareline and data stored at the second storage node with data presented on second compareline and selectively pulls down a precharged matchline in response.\nAccording to an additional embodiment of the principles of the present invention, a multiple matchline content addressable memory is disclosed. This content addressable memory includes a plurality of conductors and a plurality of memory cells, each cell comprising a storage element and a transistor for selectively coupling the storage element to a corresponding one of the plurality of conductors. A plurality of comparator circuits are included, each for comparing a bit of data stored in the storage element of a selected one of the memory cells with a bit of a comparand presented on a corresponding one of the plurality of conductors and selectively controlling an associated one of a plurality of matchlines in response.\nThe principles of the present invention are also embodied in a telecommunications subsystem including a plurality of input ports for a receiving stream of information and a plurality of output ports, a selected one of the output ports transmitting the stream of information received by the input port to a next node in a telecommunications network. A content addressable memory stores a translation table, the translation table containing information for selecting the selected one of the output ports and includes a plurality of conductors, a plurality of memory cells and comparator circuitry for comparing data presented on a selected one of the plurality of conductors with data stored on storage elements of a corresponding pair of memory cells. In response to the comparison, the comparator circuitry selectively pulls down a voltage on an associated matchline.\nThe principles of the present invention allow for the construction of and operation of content addressable memory cells with multiple comparelines and/or multiple matchlines. Moreover, in some embodiments, two pairs of bitlines are provided such that one set of bitlines can be used to access the storage elements of the memory cells while the second set of bitlines is being precharged. They are advantageously used for simultaneous updates of the stored data base while match operations are being performed. At the same time, multiple comparelines and matchlines allow for multiple comparisons to be made against multiple comparand bits simultaneously."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of memory interface design and, more particularly, to cache design in a graphics system.\n2. Description of the Related Art\nA computer system typically relies upon its graphics system for producing visual output on the computer screen or display device. Early graphics systems were only responsible for taking what the processor produced as output and displaying that output on the screen. In essence, they acted as simple translators or interfaces. Modern graphics systems, however, incorporate graphics processors with a great deal of processing power. They now act more like coprocessors rather than simple translators. This change is due to the recent increase in both the complexity and amount of data being sent to the display device. For example, modern computer displays have many more pixels, greater color depth, and are able to display images that are more complex with higher refresh rates than earlier models. Similarly, the images displayed are now more complex and may involve advanced techniques such as anti-aliasing and texture mapping.\nSince graphics systems typically perform only a limited set of functions, they may be customized and therefore far more efficient at graphics operations than the computer's general-purpose central processor. While early graphics systems were limited to performing two-dimensional (2D) graphics, their functionality has increased to support three-dimensional (3D) wire-frame graphics, 3D solids, and now includes support for three-dimensional (3D) graphics with textures and special effects such as advanced shading, fogging, alpha-blending, and specular highlighting.\nWith each new generation of graphics system, there is more image data to process and less time in which to process it. This consistent increase in data and data rates places additional burden on the memory systems that form an integral part of the graphics system. One example of a memory sub-system defining the upper limit of overall system performance may be the texture buffer of a graphics system. Certain graphics applications such as 3D modeling, virtual reality viewers, and video games may call for the application of an image to a geometric primitive in lieu of a procedurally generated pattern, gradient or solid color. In these applications, geometric primitives carry additional mapping data (e.g., a UV, or UVQ map), which describes how the non-procedural data is to be applied to the primitive. To implement this type of function, a graphics system may employ a texture buffer to store two dimensional image data representative of texture patterns, “environment” maps, “bump” maps, and other types of non-procedural data.\nDuring the rendering process, the mapping data associated with a primitive may be used to interpolate texture map addresses for each pixel in the primitive. The texture map addresses may then be used to retrieve the portion of non-procedural image data in the texture buffer to be applied to the primitive. In some cases (e.g., photo-realistic rendering) a fetch from the texture buffer may result in a neighborhood or tile of texture pixels or texels to be retrieved from the texture buffer and spatially filtered to produce a single texel. In these cases, four or more texels may be retrieved for each displayed pixel, placing a high level of demand on the texture buffer. Thus, poor performance of the texture buffer is capable of affecting a cascading degradation through the graphics system, stalling the render pipeline, and increasing the render or refresh times of displayed images.\nIn other words, accesses of graphics data, such as texture map data, must be performed very quickly. Therefore, one goal of a graphics system is to improve the speed and efficiency of memory accesses of texture maps from a texture memory. One common method is to use a texture memory cache to improve the speed of accesses of texture maps from the texture memory. The design of texture memory systems, including texture cache memory systems, plays a significant role in the implementation of new generation graphics systems.\nTexture mapping hardware often needs to process multiple pixels in the same cycle in order to maximize performance. While these pixels are independent, they typically exhibit some degree of spatial coherence. A high performance texture cache would take advantage of this coherence without imposing ordering or synchronization restrictions between pixels. Texture mapping is generally a read-only operation. Consequently, latency is typically not a factor when considering proper operations, but bandwidth does affect performance. In contrast, the performance of microprocessor instruction and data caches is typically affected by latency, not bandwidth.\nIn a texture mapping graphics system that processes multiple pixels, typically in multiple parallel pipelines, it is common for neighboring pixels to request overlapping texture data. Traditionally, cache memory in such a system requires one port per texel per pipeline, which is generally expensive to implement, or may suffer performance loss from having too few read-ports. Alternately, multiple cache memories may be built into the texture mapping graphics system. This solution is typically expensive to implement as well.\nTherefore, new texture cache memory systems and methods are desired to optimize texture caching by reducing the number of cache read-ports necessary to support the simultaneous reading of texels for multiple pixels, taking advantage of the spatial locality of the texel accesses. More generally, improved cache memory systems are desired in various different applications, including graphics applications."}
{"text": "This invention relates generally to rotating machinery and more particularly to apparatus for balancing rotors.\nGas turbine engines typically include several rotor stages, each having a rotor disk carrying an array of airfoils, i.e., compressor or turbine blades. Turbine rotors must be balanced to prevent damage and excessive loads on bearings and supporting structures, as well as efficiency losses caused by loss of clearance between the airfoils and the surrounding structure (caused by, e.g., shroud rubs).\nDespite efforts to first balance their constituent components, turbine rotors still require dynamic balancing following assembly. For this purpose, it is desirable to use balance weights that can be re-positioned to redistribute the mass of the rotor as needed and allow the system unbalance to be fine-tuned to meet precise requirements. Separable balance weights are a common practice in larger gas turbine engines. These include bolts, washers, nuts and other fasteners of varying sizes.\nIn some gas turbine rotors, notably those in smaller engines, CURVIC couplings and friction joints are assembled using a single bolt or a group of bolts (referred to as a “tie rod” or “tie bolts”) spanning the length of the assembly. A tie bolt configuration weighs less than a conventional bolted joint, but the absence of bolt holes eliminates convenient features on the rotor disk which could otherwise be used to attach separable balance weights. Accordingly, the current state of the art for smaller turbine engines is to balance the assembly by selectively machining a sacrificial surface on the rotating part. Material is removed at the location of peak unbalance to redistribute the mass of the rotor about the axis of rotation. This process is irreversible and risks damaging a component such as an integrally-bladed rotor or “blisk”, which is both safety-critical and expensive."}
{"text": "Hydronic radiant floor, wall, and/or ceiling thermal control systems are typically used for heating a space, such as a room in a dwelling or commercial building, usually for human and creature comfort. Typical hydronic heating systems require a supply of hot water, or other fluid, from a boiler, for example, and valves for controlling the quantity of the water from the supply that is fed to heating loops, which include tubing and/or heating elements. These heating loops are in thermal contact with the space or rooms to be heated or cooled. It is often necessary to control the temperature of the water in the heating loops. For example, if the supply water temperature is set to about 180° F. (80° C.) for laundry, it must be modulated down to about 100° F. or 40° C., (or lower) for radiant systems. On the other hand, chillers can be used to supply cold water during the summer months for cooling.\nA suitable system for reducing and controlling the supply water temperature is described in U.S. Pat. No. 5,119,988, issued Jun. 9, 1992, entitled “Hydronic Heating Water Temperature Control System”, to Joachim Fiedrich, the inventor herein. In that patent, a three-way, modulated diverting or by-pass valve is provided, in the return line to the boiler, for diverting some of the cooler return water to the hot supply water to reduce the temperature of the supply water feeding the heating loop supply header. This is sometimes called temperature dilution and the diverting valve is modulated by a feedback signal derived from the diluted water temperature.\nA number of approaches are available for distributing the tubing in the space to be heated to form the heating loops. In one such approach, the tubing is set in the concrete flooring. In other cases, the heat loop tubing is installed between the floor or ceiling joists using metal radiation plates.\nA preferable approach for forming the heating loops in floors relies on modular panel heating elements. Some examples are described in U.S. Pat. No. 5,292,065, issued Mar. 8, 1994, entitled “Radiant Floor and Wall Hydronic Heating Systems”, to Joachim Fiedrich, the inventor herein. The panel elements include integral metal radiation plates or sheets that are attached to two spaced apart boards, which cooperate to hold the tubing in intimate thermal contact with the radiation plate, so that the plate is heated by conduction of heat from the tubing. The plate then provides a surface that radiates heat into the room. Thermal conduction from the tubing to the plate and mechanical attachment of the tubing to the panel can also be ensured by using a resilient, thermally-conductive filler material as described in U.S. Pat. No. 5,579,996, issued Dec. 3, 1996, entitled “Radiant Floor and Wall Hydronic Heating Systems”, also to Joachim Fiedrich, the inventor herein.\nThese hydronic thermal control systems can also be used for cooling. The cooling is accomplished by feeding cool water or fluid to the tubing to reduce the temperature of the radiation plate in the modular panel, to below room temperature. As a result, heat is radiated from the room to cool water in the tubing. This heats the water slightly, and the water is then fed to a heat exchanger or chiller, for example, where it gives up the heat and is fed back to the panels.\nMore recently, the instant inventor described a system of installing the tubing in sheets of gypsum or cement wallboard as described in U.S. Pat. Appl. Publ. No. 2004/0026525 A1, entitled “In radiant wall and ceiling hydronic room heating or cooling systems, using tubing that is fed hot or cold water, the tubing is embedded in gypsum or cement wallboard in intimate thermal contact therewith so that the wallboard heats or cools the room”, which is incorporated herein in its entirety by this reference.\nOften, these modular panel systems use a number of different types of panels to create the continuous tracks required to hold the tubing of the radiant loops. Most commonly, straight, lateral run track panels are connected end-to-end to provide tubing tracks that extend laterally across the room floor to be heated or cooled. At the end of the tracks on each of these lateral run panels, “U” turn or return track panels are usually used. These return track panels comprise arcuate tracks that allow the tubing to be routed between successive tracks in the lateral run track panels by laying the tubing through the 180 degree arc of the return track. Using the combination of the straight tracks of the lateral run track panels and the return tracks of the return track panels, large serpentine radiant tubing loops can be created in the floors of rooms or other spaces.\nIn order to complete the radiant heating/cooling loops, connections must be further made between the tubing loops and the manifold, circulating pump, and/or injection valve control assemblies that are located, for example, in a closet or other area near or in the room or space to be heated or cooled. This routing between the typically serpentine layout of the tubing in the floor, for example, and the manifold, pump, and injection valves of the control assemblies can be performed either in-plane and/or out-of-plane.\nIn in-plane routing, the connections are routed, at least in part, in the plane of the floor. Often, the long runs to the control assembly can be made in tracks constructed from the lateral run track panels. This has advantages since the tubing routed in this connection can also contribute to the heating and/or cooling of the space.\nRouting between the serpentine tubing layout on the floor and the control assembly can also be performed out-of-plane. In this case, a hole is usually drilled through the floor, for example, and then the tubing is routed between or through the floor joints to connect the serpentine layout with the control assembly."}
{"text": "Businesses often receive phone calls from potential customers wishing to get information regarding products or services offered by the business. For example, the potential customer may want to know about features and pricing of a particular product. In addition, current customers may call wishing to change a service or get technical support about a particular product. Given that a business may offer numerous services or have numerous products in numerous product lines, a customer representative fielding the phone call may have difficulty in being able to provide accurate and detailed information about any one service or product. While the customer representative may have access to a computing device or computer terminal on which he or she can search for relevant information during the phone call, the process of searching may require time during which the customer has to wait for the representative to type in appropriate keywords or search strings in response to the customer's queries. This may lead to inefficient usage of both the customer's time and the representative's time, and negatively impact the customer experience."}
{"text": "The invention relates to a screw rotor machine for a compressible working medium with a screw cam rotor provided with gudgeons at its ends and a hollow screw thread rotor provided with bearing surfaces and surrounding the screw cam rotor which are rotatably journalled in a housing for rotation around mutually sideways displaced rotation axes which are fixed relative to the housing and coaxial with the gudgeons and the bearing surfaces respectively which rotors during rotation form chambers between their screw cams and screw threads which chambers move from end to end of the rotors while changing their volume.\nIn screw rotor machines of this and the closely related type where the screw thread rotor is stationary and forms a stator there are a number of suggestions as regards the choice of profile for the cooperating screw cams and screw threads of the rotors. This is shown for example in Swedish patent specification No. 85,331 and U.S. Pat. specifications Nos. 1,892,217 and 2,553,548. Hitherto existing combinations of working profiles with associated rotor journallings and drivings have, however, mostly been complicated, difficult to manufacture, and hardly suitable for effective work at high pressures and rpm."}
{"text": "Computer security is particularly challenging because it imposes negative constraints on a computer, for example: “No buffer can be overflowed.” Arguably, making a computer correctly do what it is designed to do is hard enough. In computer security, and specifically in the area of security mitigations, the problem is to prevent a computer from doing a subset of what it is (unintentionally) designed to do. This is even harder. One strategic approach to computer security is the conversion of negative requirements to positive enforceable rules, by for example, using programming strategies that make computer programs more dependable and more resistant to subversion. However, in this approach, the same developer who is responsible for introducing code defects is asked to follow various coding rules. Inevitably, in the same way that occasionally a code defect creeps in, occasionally a developer will neglect to follow one of the coding rules."}
{"text": "This invention relates to waveguides for propagating electromagnetic energy and more particularly to energy-absorbing terminations for waveguide.\nModern communications systems require ever increasing bandwidths in order to accommodate the increasing amount of data being handled. As signal bandwidths increase, the frequencies of carriers upon which these signals are superimposed have become increasingly higher. At carrier frequencies below 100 Megahertz (MHz), transmission lines are for the most part in the form of coaxial cable or open-wire transmission lines. In the region between 100 MHz and about 8 Gigahertz (GHz), coaxial cables and waveguides both find use, with waveguides being used for those applications requiring lower loss or high power handling capability. At X-band (8.2 to 12.4 GHz) and beyond, the use of waveguides predominates over the use of coaxial cable. Waveguides are currently commercially manufactured for use at frequencies up to 140 GHz.\nComplex waveguide assemblies are used for many purposes and find special use in communications satellites. For example, a microwave solid state transmitter for a satellite may be implemented by means of a large number of solid state power sources, each producing a few watts of signal energy which is either generated in or coupled into a waveguide. The signals from these power sources are combined by a \"tree\" type of waveguide power combiner, in which the power from the signal sources is applied in pairs, equal in amplitude and quadrature in phase, from the signal sources to inputs of waveguide directional couplers. Each directional coupler has two input ports and an output port, and also has a terminated waveguide port. Each directional coupler combines the signal powers which are equal and in quadrature, from the two input waveguides and sums their power in a single output waveguide. In this fashion, the signals from eight solid-state power sources applied by eight waveguides to four directional couplers may be added in pairs to produce twice the amount of power in each of four waveguides. The four waveguides are coupled in pairs to the inputs of two further directional couplers, which combine the signal powers into two waveguides which are coupled to a final directional coupler, which combines the last pair of signals onto a single waveguide. Such an arrangement requires N-1 directional couplers to combine the signal from N input waveguides. It will be understood that when the signal from large numbers of waveguides is to be combined, many directional couplers and a substantial number of waveguide connections are required. When such an assemblage is to be used on a satellite, it is imperative that the size and weight be minimized. It is clear that conventional \"plumbing\" assembly, in which each waveguide component has individual waveguide connection flanges interconnected by suitable waveguide lengths and elbows, has excessive weight due to the interconnection waveguides and flanges, and also has excess volume because of the space between waveguide components. A more satisfactory fabrication technique for an assemblage of waveguide components is to mill the entire waveguide assembly into a solid \"monolithic\" block of metal, and insert the requisite components. It is convenient when fabricating such as assemblage to fabricate two mating halves of a block which when mated together define both the waveguide components and their interconnection waveguides.\nAmong the waveguide elements required for use with a directional coupler is a terminating load. Such loads are intended to absorb all the energy flowing through a waveguide without reflection. When the two inputs are in phase quadrature, no energy flows into the terminating load. However, for any inequality in the amplitude of the inputs, or for any phase error, energy flows into the terminating load and must be absorb to reduce reflects back to the input. In the prior art, as described for example in U.S. Pat. No. 3,904,993 issued Sept. 9, 1975, to James, the waveguide load consists of a wedge of energy absorbing or lossy material located in and extending across the full width (larger cross-sectional dimension) of the rectangular waveguide. In the James arrangement, conductive septums are located within the lossy material to aid in carrying away heat. It has been found that such prior art terminations create difficulties in assembling the two mating portions of a waveguide assemblage such as that previously described. For example, an 8-way power combiner including seven directional couplers requires seven waveguide terminations or loads. As mentioned, the prior art energy absorbing material extends across the full width of the waveguide. The energy absorbing wedges, once cut to size, can be inserted into the proper location in one of the halves in the assembly. The wedges are retained in that location by use of adhesive. When so assembled, each wedge protrudes above the half of the waveguide into which it is inserted by a full half width of the waveguide. Thereafter, the mating half of the assemblage is placed on the half containing the absorbent wedges. It has been found to be extremely difficult to simultaneously fit all of the wedges into the mating halves of the waveguide. This difficulty exists even when the absorbent wedges are formed from a material which is resilient, such as an elastic foam. This problem is exacerbated if those portions of the absorbent material which are intended to be in contact with the walls of the mating waveguide are coated with adhesive. Even if assembly can be accomplished, the adhesive once cured prevents disassembly of the two halves of the waveguide assemblage and, if the assembly is forcibly disassembled, tears the absorbent material in a manner making reassembly difficult.\nA waveguide termination arrangement is desired which is amenable to convenient assembly and disassembly of the mating halves of a waveguide assemblage."}
{"text": "Films produced in foreign countries are sometimes dubbed into the mother tongue or added with the native language subtitles before they are released. Bilingual TV broadcasting systems have been developed, in which the user can select the desired language at his or her TV set.\nThe user can select either of the two languages but cannot select any other language.\nThe market of systems using recording mediums (films, disks, magnetic tape, etc.) has expanded world-wide. Demand Has therefore grown for systems capable of reproducing recorded information in more than one, or even two, languages."}
{"text": "The invention concerns cellular radio systems, such as mobile phone systems and PMR systems (PMR=private mobile radio), i.e. cellular radio networks in which the area covered by the network is divided into radio cells. The subscribers roaming in the area of these radio cells communicate with the system when they are within the coverage area of the base stations located in the radio cells.\nIn a system like this, a subscriber station, such as a mobile phones, located in a radio cell listens or tunes the receiver unit of its radio to the frequency of a signalling channel of a base station located in the radio cell. A signalling channel is used for exchanging messages concerning the connection establishment between a subscriber station and a base station. On the basis of connection establishment, the system controller of the cellular radio network assigns the subscriber stations to actual traffic channels, in which the actual transfer of information, i.e. speech or data, takes place. Traffic between a base station and subscriber stations takes place on radio channels that are implemented such that downlink traffic (from a base station to a subscriber station) takes place at a first frequency, and uplink traffic (from a subscriber station to a base station), at a second frequency. In FDMA systems (FDMA=Frequency Division Multiplexing Access) a channel thus means a pair of channels used for one conversation. In a TDMA system (TDMA=Time Division Multiplexing Access) a channel means a carrier and its time slot.\nIt is typical of channels provided between base stations and subscriber stations that there is a limited number of channels available since there are few radio frequencies and/or time slots. It is thus typical of cellular radio systems that one--possibly the only--signalling channel of a base station can be converted to a traffic channel, e.g. a speech or data channel, when a traffic channel is needed for a call between a subscriber station and a base station and the actual traffic channels of the base station are busy. Typically, this kind of procedure is followed in areas where the traffic and subscriber density is low and where the number of radio channels is small; the use of a signalling channel as a traffic channel significantly improves the transmission capacity of the system. When a signalling channel that can be converted to a traffic channel (i.e. a non-dedicated control channel) is used as a traffic channel, signalling messages cannot be transmitted on said channel unless the channel has an inband signalling channel, such as an associated signalling channel, whereby it is possible to transmit short signalling messages in the middle of speech or data to the parties listening to the channel (base station and subscriber station). However, it is not possible to set up new calls on the associated signalling channel nor to transmit data messages on it since the transmission capacity of said channel is very limited. Typically all the channels of a base station can function as signalling channels or traffic channels when necessary and that if all the channels of the base station function as traffic channels at the same time, any channel that is released from use as a traffic channel can be converted to a new signalling channel for traffic between the base station and the subscriber stations roaming within its coverage area. This is optimal in respect of formation of a new signalling channel: a new signalling channel is formed without delay. When a channel like this is used as a signalling channel, signalling messages concerning all the services of a mobile phone system can be transmitted thereon. It should be noted that a radio system may simultaneously comprise base stations with dedicated control channels only and base stations with non-dedicated control channels.\nIt is typical of cellular radio systems, particularly of PMR systems, that the radio cells or the coverage areas of the base stations intersect or overlap, whereby a subscriber station can select--by measuring the signals transmitted on the signalling channels by the base stations--the base station whose signalling channel (control channel) it will start to listen to. The coverage areas of the base stations may overlap heavily. The coverage areas overlap especially in low-lying and flat areas when great transmission power is used, e.g. in PMR systems.\nDue to the overlapping of the coverage areas of the base stations in the prior art, when all the channels of a base station are converted to traffic channels, the subscriber stations which have listened to a base station channel that has functioned as a signalling channel but has been converted to a traffic channel move to a signalling channel of an adjacent base station. Since in normal mobile phone traffic, tele-communication links, such as speech or data links, are established from and to subscriber stations and since the adjacent base stations have few free traffic channels or none at all when the traffic load is normal/optimal, all the traffic channels and non-dedicated control channels of the adjacent base stations are quickly allocated for use as traffic channels. Those subscriber stations that do not find a signalling channel to which they could be tuned or a traffic channel on which they could communicate with the base station start to tune their radio units to the frequencies of the different base stations, attempting to find a free signalling channel and a traffic channel if necessary and simultaneously attempting to register in the network and send the parameters for establishing a telecommunication link with several base stations in succession yet always failing since the other base stations do not have free channels either. The subscriber stations thus roam between the coverage areas of different base stations, causing instability and unnecessary signalling load in the network. Further, the fact that the subscriber stations roam to listen to the signalling channels of base stations that are not their original base stations may lead to a situation where not a single subscriber station is listening to the signalling channel of the original base station when a signalling channel or even a traffic channel of said original base station becomes available, whereby the network has a free signalling channel and traffic channel although at the same time the subscriber stations roaming in the area have too few traffic channels in their use. The network capacity is thus divided unevenly and resources are wasted. The subscriber stations detect this as a decrease in the quality of service: it is difficult for them to find a signalling channel at a base station and to establish telecommunication links, and their call set-up attempts are interrupted since often the signalling channels on which a subscriber station attempts to set up a call, i.e. to send messages necessary in call set-up, is allocated to another subscriber station for use as a traffic channel."}
{"text": "1. Field of the Invention\nThe present invention relates to a multiplexing device which multiplexes a plurality of pieces of data applied.\n2. Description of the Related Art\nVarious standards related to multimedia communication have been specified so far. ITU-T Recommendation H.324, for example, defines a system and a terminal device for multimedia communication using an analog channel. Also specified are Standard H.324/M for a mobile channel switched network and Standard 3G-324M for an advanced mobile communication network.\nIn these standards, Standard H.223 is specified as a system for multiplexing control data, voice coding data and image coding data into a single stream.\nStandard H.223 provides a logical channel for each of voice data, image data and control data to enable QoS (Quality of Service) control for each logical channel.\nFor H.324/M and 3G-324M, the system which has improved transmission error resistance of Standard H.223 is specified by Standard H.223 Annex A,B,C.\nBased on such standards, systems and the like for multiplexing voice data, image data and control data have been proposed (see e.g. Japanese Patent Laying-Open No. 2002-152730) and developed.\nCommon multiplexing device according to conventional art which can be expected from ITU-T Recommendation H.223 will be described.\nFIG. 6 is a block diagram showing a structure of a common multiplexing device. With reference to FIG. 6, the common multiplexing device includes a multiplexing parameter control unit 103, a control data adaptation unit 104, a voice data adaptation unit 105, an image data adaptation unit 106 and a multiplexing layer unit 107.\nThe multiplexing device shown in FIG. 6 is a device which multiplexes coded image and voice into a single stream. To the present multiplexing device, applied other than image data 102 which is a coded image and voice data 101 which is coded voice are control data 100 as control information related to a coding system or a multiplexing system of voice and images.\nThe control data 100 is input to the control data adaptation unit 104, the voice data 101 is input to the voice data adaptation unit 105 and the image data 102 is input to the image data adaptation unit 106.\nThe multiplexing parameter control unit 103 assigns each of the control data adaptation unit 104, the voice data adaptation unit 105 and the image data adaptation unit 106 its maximum AL-SDU (Adaptation Layer-Service Data Unit) size.\nAL-SDU is an SDU of an adaptation layer and a maximum AL-SDU size is a parameter which specifies a maximum size of an AL-SDU. Assume that a maximum AL-SDU size of the control data adaptation unit 104 is a maximum AL1-SDU size, a maximum AL-SDU size of the voice data adaptation unit 105 is a maximum AL2-SDU size and a maximum AL-SDU size of the image data adaptation unit 106 is a maximum AL3-SDU size.\nThe multiplexing parameter control unit 103 assigns the multiplexing layer unit 107 a maximum MUX-PDU (Multiplex-Protocol Data Unit) size.\nMUX-PDU is a PDU in a multiplexed layer and a maximum MUX-PDU size is a parameter which specifies a maximum size of an MUX-PDU.\nThe control data adaptation unit 104, the voice data adaptation unit 105 and the image data adaptation unit 106 provide a function of an adaptation layer according to properties of each data.\nThe control data adaptation unit 104 generates an AL-PDU (AL1-PDU) which stores applied control data 100 and sends the same to the multiplexing layer unit 107. At this time, when the size of the input control data 100 is equal to or smaller than the maximum AL1-SDU size, the control data adaptation unit 104 takes the control data 100 as an AL-SDU (AL1-SDU) without processing. When the size of the control data 100 is larger than the maximum AL1-SDU size, the control data adaptation unit 104 segments the control data 100 on a maximum AL1-SDU size basis to take each segmented data as the AL1-SDU. As to the control data, no additional information exists, so that the AL1-SDU is as it is taken as the AL1-PDU.\nThe voice data adaptation unit 105 generates an AL-PDU (AL2-PDU) which stores applied voice data 101 and sends the same to the multiplexing layer unit 107. At this time, when the size of the input voice data 101 is equal to or smaller than the maximum AL2-SDU size, the voice data adaptation unit 105 takes the voice data 101 as an AL-SDU (AL2-SDU) without processing. When the size of the voice data 101 is larger than the maximum AL2-SDU size, the voice data adaptation unit 105 segments the voice data 101 on a maximum AL2-SDU size basis to take each segmented data as the AL2-SDU. As to the voice data, a sequence number and a CRC (Cyclic Redundancy Check) code are attached to the AL2-SDU as required to have the AL2-PDU. Sequence number is a number indicative of the order of data and a CRC code is a code for use in detecting a bit error.\nThe image data adaptation unit 106 generates an AL-PDU (AL3-PDU) which stores applied image data 102 and sends the same to the multiplexing layer unit 107. At this time, when the size of the input image data 102 is equal to or smaller than the maximum AL3-SDU size, the image data adaptation unit 106 takes the image data 102 as an AL-SDU (AL3-SDU) without processing. When the size of the image data 102 is larger than the maximum AL3-SDU size, the image data adaptation unit 106 segments the image data 102 on a maximum AL3-SDU size basis to take each segmented data as the AL3-SDU. As to the image data, a sequence number and a CRC code are attached to the AL3-SDU as required to have the AL3-PDU.\nThe multiplexing layer unit 107 multiplexes the AL1-PDU from the control data adaptation unit 104, the AL2-PDU from the voice data adaptation unit 105 and the AL3-PDU from the image data adaptation unit 106 and attaches an MUX-PDU header to the multiplexed data to generate an MUX-PDU and output the same as multiplexed data 108. At this time, the multiplexing layer unit 107 segments the AL-PDU as required such that the size of the MUX-PDU is within the maximum MUX-PDU size. With a plurality of multiplexing tables for use in multiplexing prepared in advance, the multiplexing layer unit 107 uses any of the multiplexing tables for the multiplexing processing.\nMUX-PDU header is formed of a PM (Packet Marker), an MC (Multiplex Code) and an HEC (Header Error Control). PM indicates whether a last byte of an AL-SDU is included or not. MC indicates a multiplexing table used. HEC is information for detecting an error in an MUX-PDU header.\nWith a common multiplexing device, a size of data which can be processed is limited due to constraints on hardware and the like. In general, however, the image data 102 has a larger size than those of the control data 100 and the voice data 101, so that the image data 102 whose size exceeds the upper limit might be applied in some cases. In such a case, the image data 102 should be segmented or abandoned. When the data is abandoned, it is impossible to properly multiplex applied data.\nIn addition, while the image data 102 is in general input on a VP (Video Packet) basis or a segmented data basis, if the multiplexing device simply segments the image data 102 irrespectively of a delimitation of a VP, the delimitation of the VP goes out of order to prevent proper decoding of multiplexed image data."}
{"text": "This invention relates to a replacement wheel bearing spindle which may be an accessory for boat trailers or the like and more particularly to a replacement bearing spindle which may be attached to a damaged bearing spindle at the wheel carrying end of a trailer axle to permit the wheel and new bearings to be mounted on the axle.\nOn a boat trailer or the like a fixed axle is carried at one end by the trailer chassis and on the other end has an integral spindle which generally carries a pair of bearings at spaced locations. The bearings have inner races fastened to the spindle and outer races on which a wheel is mounted for rotation relative to the spindle. A tire hub is fastened to the wheel, supports the trailer and rotates relative to the trailer chassis. When a bearing fails or \"burns-out\" the heat generated resulting from the friction created while the trailer is being pulled and before the failure is noticed, tends to damage the spindle by enlarging the diameter of the surface which seats the inner race of the bearing. When this occurs the trailer cannot be moved since a spare new set of bearings will not be received on the spindle seating surface and the trailer must be towed at great expense for subsequent replacement of the axle. If an axle is not available at the repair station, the trailer may be out of commission for some time.\nAlthough such failures appear to occur relatively frequently, solutions to place the trailer back into service even temporarily do not appear to have been provided by the prior art. A solution to this problem must of course be simple enough to perform in the field while the trailer is on the road and be cost effective."}
{"text": "Drill cuttings are the earth, rock and other solid materials generated during the rotary drilling of subterranean wells. The drill cuttings are removed from beneath the drill bit by a stream of drilling fluid that suspends the solids and carries the solids to the surface. On the surface, the drill cuttings are separated from the drilling fluid in a drilling cuttings separator or shaker and the drill cuttings are collected at the drill site for subsequent treatment.\nTraditional oleaginous drilling fluids, also known as oil-based drilling fluids or invert emulsion drilling fluids (if they contain an internal non-oleaginous phase), may be harmful to marine life due to the presence of aromatic hydrocarbons in the diesel fuel or other similar petroleum fractions used as the continuous phase. The development of low-toxicity mineral oil-based drilling fluids—with very low fractions of aromatic compounds—allayed much, but not all, of the concern over acute toxicity effects on marine flora and fauna. However, discharge of mud-laden cuttings still produces a mound of cuttings on the ocean floor that may smother any marine life that resides on the seabed.\nDevelopment of synthetic-based drilling fluids as alternatives to conventional oil-based drilling fluids in offshore operations was precipitated by residual toxicity and biodegradability concerns. These developments focused on the fate and effects of oil-coated drilled cuttings discharged into the sea, as well as worker safety. For onshore applications, cuttings disposal is also of importance. However, since the drilled cuttings are disposed of on land, the environmental issues focus primarily on subsequent usability of the land and contamination of ground water. Although the advent of synthetic-based fluids has greatly improved the environmental acceptability of non-aqueous drilling fluids both offshore and onshore, current synthetic-based fluid formulations still present problems for direct land treatment of oil-coated cuttings resulting from onshore operations. The concerns with pollution of soil and groundwater by synthetic based fluids and oil-based fluids have led to increasingly strict government regulations.\nOily drill cuttings can have severe impacts on their receiving environment and should be cleaned or treated to minimize their environmental impact and the operator's long term liability. The primary purpose of each of these methods is to somehow destroy or remove the drilling fluid residue from the earth solids. In addition to the above mentioned method of land treatment (spreading and farming), there is a litany of other ways to treat oil-coated cuttings from drilling operations. These include landfill disposal; bio-remediation; stabilization/solidification (briquetting, fixation with silicates or fly ash); extraction or washing (oil, detergents, and solvents); and thermal treatment (incineration and distillation, including thermal desorption and hammer mill). The treatment of drill cuttings is the subject of a number of patent applications and literature disclosures that include U.S. Pat. Nos. 6,187,581; 6,020,185; 5,720,130; 6,153,017; 5,120,160; 5,545,801; 4,696,353; 4,725,362; 4,942,929; 5,132,025. These patents describe various methods of treating oily drill cuttings including incineration; reinjection of the slurrified cuttings into another subterranean formation; chemical washing and landfill disposal; and other methods. As noted above, the primary purpose of each of these methods is to somehow destroy or remove the drilling fluid's residue from the earth solids.\nDespite considerable research conducted in the area of drill cuttings disposal, there remains an unmet need for a clean, inexpensive and environmentally friendly drilling fluid and method of treating the drill cuttings such that they produce an end product that may have a beneficial use."}
{"text": "In general terms voice communication can be said to involve the transmission of a near-end speech signal to a far-end or distant user, where a speech enhancement problem consists in the estimation of a relatively clean speech signal from a captured noisy signal. There are a number of single-microphone configurations which allow for improvements when considering the suppression of noise.\nUse of two distinct microphones to simultaneously capture a sound field allows for a possible usage of spatial information and characteristics of the sound source(s) from which a sound field captured by the microphones originates. These characteristics may relate to the relative placement of the microphones on a mobile communication device as well as the design and usage of the communication device. A proper estimation of the noise characteristics forms a basis for an efficient use of noise suppression algorithms, such as e.g. algorithms which are based on spectral subtraction, which is commonly used in this particular technical field.\nDifferent methods for executing dual-microphone noise suppression have been suggested based on the assumption that the signals received by the microphones have a relatively similar power level for the near-end signal generated by the user of the communication device.\nIn WO 2007/059255 noise suppression is performed by generating a ratio of power difference and sum signals from input signals captured by two microphones, after which the input signals are being processed such as to suppress the estimated noise from one of the two input signals.\nA drawback with WO 2007/059255, which is relying on the assumption of small or even no gain difference between signals captured by a microphone pair is that, in practice, dual-microphones mounted side-by-side on mobile devices will present an arbitrary gain difference. This difference is both inherent to the high variation of the manufactured microphone gains and to the variation in the near-field signal received levels with small changes in the position of the mobile device relative to the speaker's mouth, when the device is used in handheld mode.\nOther methods, such as e.g. the one presented in US 2007/0154031 exploit the level differences between received microphone signals to discriminate speech and noise in the time-frequency domain and to suppress the noise accordingly.\nHowever, while the use of a microphone for capturing noise, typically referred to as a reference microphone, in conjunction with a microphone used for capturing basically speech, typically referred to as a primary microphone, and the exploitation of a resulting signal level difference at the two microphones can allow for a fairly good detection of the speech and noise signals in the time-frequency domain, noise suppression based on a masking approach, such as the one described in US 2007/0154031 normally results in a high distortion of the extracted speech signal and introduces also often musical noise.\nA spectral subtraction based method applicable for dual-microphone noise suppression has been suggested in WO2000/062579, where spectral processors are used for producing separate noise reduced and noise estimated signals.\nSpectral subtraction techniques, such as the one described in WO2000/062579, have generally proven to be relatively robust to speech cancellation and to provide a relatively good suppression of stationary noise. The filtering process which is normally used in association with spectral subtraction usually relies on estimates of the spectrum of the noise and the spectrum of the noisy speech. The noise spectrum is preferably estimated during speech pauses and is based on the estimation of the stationary part of the noise only. Many background noise environments, such as e.g. restaurants, airports, streets and other public places, are however characterized by the presence of a high level of non-stationary noise which is not taken into consideration in known implementations, which are based on spectral subtraction techniques, and hence when applying these techniques the non-stationary noise component remains unfiltered in the signal transmitted to the far-end user of the communication link."}
{"text": "In a warehouse or other location from which ordered goods are shipped to an ordering party, such as a fulfillment center, items of stock progress through a number of stages after being ordered from a vendor and before being shipped to a customer. For example, an inbound process for items of stock may include delivery of the items of stock at a warehouse, a sort process in which the items of stock are sorted by a first worker into a number of categories, a receive process in which the items of stock are entered into inventory by a second worker individually scanning each item with a portable electronic device and a stow process in which the items of stock are placed by a third worker in an appropriate storage location until being ordered by a customer. An outbound process may similarly include a number of processes. Each of these processes takes time and effort."}
{"text": "Methane comprising fluids can be derived from a number of sources, such as natural gas or petroleum reservoirs, aerobic or anaerobic digestion of biological material, or from a synthetic source such as a Fischer-Tropsch process.\nAerobic or anaerobic digestion of biological material can have many forms. It may, for example, concern aerobic or anaerobic digestion of dairy waste. Aerobic or anaerobic digestion of biological material results in a methane comprising fluid which may also be referred to as bio-methane, bio-gas or bio-fuel.\nThe Fischer-Tropsch process can be used for the conversion of synthesis gas (from hydrocarbonaceous feed stocks) into liquid and/or solid hydrocarbons. Generally, the feed stock (e.g. natural gas, associated gas and/or coal-bed methane, heavy and/or residual oil fractions, coal, biomass) is converted in a first step into a mixture of hydrogen and carbon monoxide (this mixture is often referred to as synthesis gas or syngas).\nIn this text the term natural gas is used to refer to methane comprising streams originating from any source.\nThe term methane comprising fluid is used in this text to refer to fluids comprising at least 40 m/m % methane. The term in particular relates to natural gas fluids. Similarly, the terms methane comprising gas, methane comprising liquid, and methane comprising slush of liquid and solids, refer to gas, liquid, and slush of liquid and solids comprising at least 40 m/m % methane. A methane comprising fluid may comprise up to 100 m/m % methane, especially up to 99.9 m/m % methane, it may comprise up to 99 m/m % methane, it may comprise up to 90 m/m % methane. In the context of this invention a slush of liquid and solids is considered a fluid. The term methane comprising gas in particular relates to natural gas. The terms methane comprising liquid, and methane comprising slush of liquid and solids, in particular relate to LNG.\nNatural gas is a useful fuel source, as well as a source of various hydrocarbon compounds. It is often desirable to liquefy natural gas in a liquefied natural gas (LNG) plant at or near the source of a natural gas stream to enable compact storage of the natural gas and/or efficient transport of the natural gas over long distances. Natural gas can be more easily stored and transported in a liquid form than in a gaseous form because it occupies a smaller volume.\nLiquefied natural gas plants are well known in the field and comprise the following processing steps:                optionally treating the methane comprising stream by removing impurities in a treating stage, such as water, acid gases, mercury,        optionally removing natural gas liquids from the methane comprising stream in a NGL stage, such as ethane, propane, butane and heavier components,        cooling the methane comprising stream in one or more cooling stages, for instance a pre-cooling stage and a main cooling stage, and        optionally flashing the methane comprising stream in an end-flash stage and,        optionally, storing the liquefied natural gas in a storage tank.        \nNatural gas is odorless. It is known to odorize natural gas so that people can easily detect very low concentrations of natural gas in air.\nWhen transporting, storing, or handling a methane comprising liquid or slush, normally boil off gas is created due to heat ingress. Thus, there is a need for an odorized methane comprising liquid or slush so that low concentrations of boil off gas in air can be easily detected by people.\nOdorants for a methane comprising liquid or slush, however, need to fulfill a number of requirements. The odorant needs to have a good solubility in a methane comprising liquid or slush at very low temperatures, even at temperatures around −162° C., or even −185° C. Furthermore, the odorant needs to remain in the methane comprising fluid when it transfers from the liquid phase to the gas phase. Additionally, the odorant needs to have a sufficiently low odor threshold value."}
{"text": "This invention relates to a piezoelectric monolithic crystal device, filters, resonators or the like, operating in the trapped energy mode and having a capacitor plated on the piezoelectric wafer within energy coupling distance of the main electrodes and it is an object of the invention to provide an improved device of this nature.\nPiezoelectric monolithic crystal devices of the nature here involved are known to the art, as for example in the U.S. Pat. No. 3,898,489, Aug. 5, 1975 to John J. Grady and Theodore E. Lind and assigned to the same assignee as the present invention. As disclosed in the Grady et al patent electrodes are plated on the major surfaces of a piezoelectric crystal wafer, for example, AT cut quartz. The dimensions of the electrodes and those of the crystal wafer were selected, together with the frequency, so that the device whether acting as a filter, resonator, or the like, was operating in the trapped energy mode.\nConducting tabs connected to the electrodes on the appropriate surfaces of the crystal wafer were provided, and on the opposite sides of the crystal wafer from the conducting tabs there were provided pseudo tabs opposite the conducting tabs but terminating short of the actual electrodes for trapping spurious modes propagating within the crystal structure and conducting them to the edges of the crystal wafer for dissipation. Such spurious modes, accordingly, did not appear in the output response.\nConnecting capacitors across the input and output electrodes of a monolithic crystal filter element for improving filter response also is well-known. In this respect, reference is made to the following U.S. Pat. Nos.: 3,947,784, Mar. 30, 1976 Arvanitis et al, 3,716,808, Feb. 13, 1973 Malinowski et al and 3,633,134, Jan. 4, 1972 Barrows et al all assigned to the assignee of the present invention. The capacitors illustrated in the patents referred to are not capacitors whose electrodes are plated onto the piezoelectric cyrstal surface but are separate, discrete, capacitors connected to the terminals as shown. However, it is also known to plate capacitor electrodes on the major surfaces of the piezoelectric element and connecting the capacitor electrodes to the input and output terminals. In such structures, though, the plated capacitor has been disposed far enough away from the main electrodes so that any frequency modes that were propagated in the piezoelectric material by virtue of the exitation of the capacitor did not couple to the main electrodes and thus did not produce any undesirable responses in the filter output. Thus such structures require relatively large piezoelectric crystal wafers in order to provide sufficient distance between the plated capacitor and the plated main electrodes to avoid any of the indicated coupling. Accordingly, such piezoelectric monolithic devices resulted in relatively large structures that were expensive in and of themselves and increased the cost by virtue of requiring larger packages. Moreover, it is essential that the device and the resulting package be made as small as possible."}
{"text": "More specifically, the invention concerns the production of a circuit comprising MIS transistors and capacitors, especially for amplifier circuits, switched capacity filtering circuits, analog/digital or digital/analog converter circuits and dynamic DRAM memories With these types of circuits, the capacitors are situated on the electric or field oxide insulation zones of the integrated circuit.\nIn one known analog die, the capacitors are formed by a first film of doped polycrystalline silicon noted poly-1, a thermic oxidation of this silicon film or a CVD oxide deposit and a second film of polycrystalline silicon deposited on the thermic oxide and noted poly-2. In the first polycrystalline silicon film, the gate of the transistors of the integrated circuit is formed. One example of known analog circuits of this type is described in the document FR-A-2 601 817 filed in the name of the Applicant.\nIn addition, in the document IEEE 1989--Custom Integrated Circuits Conference and entitled \"A Submicron Analog, CMOS Technology\" by R. W. Gregor and al., 18.5.1-18.5.4, instructions are contained for the embodiment of high value resistors in the poly-1.\nThe poly-2 used for the embodiment of the capacitors introduces a further level in the structure of the integrated circuits, thus significantly increasing the height of these circuits, which renders it more difficult to embody the stages for levelling (known under the term \"planarization\") by the dielectric materials used between the various conductive films. This also makes it difficult to reduce the dimensions of the integrated circuits. In fact, the smallest pattern of these circuits, generally constituted by the gate of the transistors, is 1 micrometer.\nMoreover, the etching of the poly-2 is delicate, as it is formed from the etching residue at the foot of the poly-1 owing to the abrupt action constituted by the gates of the transistors situated between the lateral electric insulation zones. So as to avoid this residue, a chemical isotropic etching needs to be used, thus resulting in a large dimensional loss on the poly-2.\nSo as to reduce this dimensional loss, generally a slight thickness of the poly-2 is selected, thus resulting in the embodiment of high value electric resistors for the upper electrodes of the capacitors limiting the operating speed of these circuits.\nIn one modified method, it is possible to embody the poly-1 of the capacitors before the poly-1 of the gates, the latter then being embodied in the poly-2. In this case, the difficulty of the etching the capacitive elements is to be added to the etching of the gates, the etching residue forming at the foot of the poly-1 capacitors. Now, the etching of the gates is a crucial operation for embodying a submicronic analog die.\nIn known logic (or digital) dies, on the gates of the transistors and on their source and drain zones, a silicide film of a refractory metal with low resistivity is formed via the reaction of the silicon and a refractory metal. This technique is known as the self-aligned silicide technique and is described in detail in the document IEEE Journal of Solid-State Circuits, vol. SC-20, No 1 February 1985 and entitled \"Development of the self-aligned titanium silicide process for VLSI applications\" by M. E. Alperin and al., pp. 61-69.\nThis technique has a certain number of advantages, such as the reduction of the resistance of the gate and the interconnections of these polycrystalline silicon gates (from 25 ohms to 3 ohms/square), thus improving the operating speed of the integrated circuits. In addition, it makes it possible to reduce the resistance of the source and drain zones by embodying silicide bond pads making it then possible to reduce the number of electric contacts on these active zones resulting in an integration density gain.\nThis silicide film only ensures connections either between the gates or the drains or between the sources, but it does not allow for the embodiment of interconnections between the sources and the drains or between the gates and the active zones (source or drain). In fact, these interconnections are effected after depositing a insulating (dielectric) film on the surface of the circuit and then opening contact openings in this insulating film in a metallic film etched according to the desired patterns; they constitute the first level of interconnections.\nThis silicide technology is incompatible with known analog dies, as they are unable to correctly oxidize a silicide so as to form silicide/oxide/poly capacitors. Furthermore, if the silicide is formed after the poly-2 is etched, the capacitors would be short-circuited by the silicide.\nIn addition, there currently exists the technique of producing local interconnection lines (LIL), which is an extension of the silicide technique described above. This technique is described in detail in the document J. Vac. Sci. Technol. B(6), November/December 1988 and entitled \"TiSi.sub.2 strap formation by Ti-amorphous-Si reaction\" by H. J. W. Van Houtum and al., pp. 1734-1739.\nIt consists of forming silicide interconnections both on the drains and sources of the transistors to be connected and on the thick dielectric (generally LOCOS). This makes it possible to ensure short distance links between the transistors without any interface being provided between the various films and especially being able to mount the electric contacts on the dielectric. It is then possible to reduce the surface of the sources and drains and thus increase the integration density of the circuits.\nIn addition, this technique has one significant advantage, namely that the opening of electric contact openings in the dielectric covering the transistors and capacitors is rendered much easier In fact, the contact openings are etched in a thick dielectric which is much more uniform than the one used in the analog die, thus facilitating the levelling of the circuits.\nAs mentioned previously, this technique is incompatible with known analog dies."}
{"text": "1. Field of the Invention\nThe present invention generally relates to target tracking and engagement techniques for achieving accurate launching of a submunition and, more particularly, is concerned with a target tracking system for determining the bearing to a ground object emitting acoustic energy.\n2. Description of the Prior Art\nA wide area mine (WAM) submunition is currently being developed for the U.S. military. The WAM (also termed a \"smart\" mine) submunition basically includes a noise sensor that can detect sounds, or an acoustic signature, of a moving target, a launcher tube containing a sublet, a positioning mechanism for aiming the launcher tube, and an arming and firing mechanism for launching the sublet from the launcher tube.\nThe WAM submunition is initially deployed in a target area in a laid down position on the ground. A self-righting mechanism of the WAM submunition is used to erect the launcher tube from the laid down position to an upright position prior to final arming. After erecting the launcher tube, the self-righting mechanism stabilizes the submunition in the upright position while maintaining it coupled with the ground. The positioning mechanism then aims the launcher tube toward the target so that the sublet when fired will be launched over the target.\nFor the sublet to be successfully launched over the target, tracking and engagement techniques must be available that are capable of determining the bearing to the sensed target and range of the target in the first place so that proper aiming of the launcher tube and launching of the sublet can be achieved. To be effective such tracking and engagement techniques must demonstrate a high degree of accuracy, have sufficient range and be capable of operating in a wide variety of environmental and site conditions.\nHeretofore, analog tracking techniques have been employed that use small microphone arrays for sensing the acoustic signature. These techniques have tracking ranges limited to less than 100 meters and are subject to a significant degree of tracking errors.\nDigital tracking techniques that have been employed heretofore involve cross-correlation and beam-forming operations. These techniques require more processing capability in the system and/or very large array sizes in order to achieve a satisfactory level of accuracy and range.\nConsequently, a need remains for improvements in tracking and engagement techniques for determining the bearing and range of a ground target in order to accurately aim a submunition toward the target at significant range."}
{"text": "1. Technical Field\nThe present invention relates to a method and system for data processing in general, and in particular to a method and system for providing an eviction protocol for a sparse directory. Still more particularly, the present invention relates to a method and system for evicting a cache line from a sparse directory within a non-uniform memory access computer system.\n2. Description of the Prior Art\nIt is well-known in the computer arts that greater computer system performance can be achieved by combining the processing power of several individual processors to form a multiprocessor (MP) computer system. MP computer systems can be designed with a number of different topologies, depending on the performance requirements of a particular application. A symmetric multiprocessor (SMP) configuration, for example, is one of the simpler MP computer system topologies that are commonly used, in which resources such as a system memory are shared by multiple processors. The topology name \"symmetric\" stems from the fact that all processors within an SMP computer system have symmetric access to all resources within the system.\nAlthough the SMP topology permits the use of relatively simple inter-processor communication and data sharing protocols, the SMP topology overall has a limited scalability and bandwidth, especially at the system memory level as the system scale increases. As a result, another MP computer system topology known as non-uniform memory access (NUMA) has emerged as an alternative design that addresses many of the limitations of the SMP topology, at the expense of some additional complexity.\nA typical NUMA computer system includes a number of interconnected nodes. Each node includes at least one processor and a local \"system\" memory. The NUMA topology name stems from the fact that a processor has lower access latency with respect to data stored in the system memory at its local node than with respect to data stored in the system memory at a remote node. NUMA computer systems can be further classified as either non-cache coherent or cache coherent, depending on whether or not data coherency is maintained among caches in different nodes. The NUMA topology addresses the scalability limitations of the conventional SMP topology by implementing each node within a NUMA computer system as a smaller SMP system. Thus, the shared components within each node can be optimized for use by only a few processors, while the overall system benefits from the availability of larger scale parallelism with relatively low latency.\nThere are several problems associated with the development of an eviction protocol for a sparse directory within a NUMA computer system. The first problem arises from the management of an entry that is selected to be evicted, the second problem arises from the generation of remote transactions that are used to cause the eviction to occur, and the third problem arises from the generation of a transaction to cause the actual writeback of data to the local memory. Consequently, it would be desirable to provide an improved method for evicting a cache line from a sparse directory within a NUMA computer system."}
{"text": "1. Field of the Invention\nThe present invention relates to an IPTV receiver and a method of discovering an IPTV service, and more particularly, to an IPTV receiver and a method of discovering an IPTV service using an Internet protocol.\n2. Discussion of the Related Art\nConventional TV services are provided such that a cable, terrestrial, or satellite broadcast provider transmits content created by a broadcaster through a radio communication medium such as a broadcast network and users view the broadcast content using a TV receiver that can receive signals of the communication medium.\nAs digital TV technologies have been developed and commercialized, it has become possible to provide a variety of content such as real-time broadcasts, Content on Demand (CoD), games, and news to viewers not only using the existing radio medium but also using the Internet connected to each residence.\nOne example of provision of content using the Internet is an Internet Protocol TV (IPTV) service. The IPTV service provides an information service, moving image content, broadcasts, etc., to a television using high-speed Internet.\nWhile the IPTV service is similar to general cable broadcasting or satellite broadcasting in that it provides broadcast content such as video content, the IPTV service is characterized in that it also supports bidirectional communication. The IPTV service also allows users to view a desired program at a desired time, unlike general terrestrial broadcasting, cable broadcasting, or satellite broadcasting."}
{"text": "1. Field of the Inventions\nThe present inventions relate to medical devices and, more particularly, to methods and apparatuses for dynamic spinal stabilization.\n2. Description of the Related Art\nThe human spine is a flexible weight bearing column formed from a plurality of bones called vertebrae. There are thirty three vertebrae, which can be grouped into one of five regions (cervical, thoracic, lumbar, sacral, and coccygeal). Moving down the spine, there are generally seven cervical vertebra, twelve thoracic vertebra, five lumbar vertebra, five sacral vertebra, and four coccygeal vertebra. The vertebra of the cervical, thoracic, and lumbar regions of the spine are typically separate throughout the life of an individual. In contrast, the vertebra of the sacral and coccygeal regions in an adult are fused to form two bones, the five sacral vertebra which form the sacrum and the four coccygeal vertebra which form the coccyx.\nIn general, each vertebra contains an anterior, solid segment or body and a posterior segment or arch. The arch is generally formed of two pedicles and two laminae, supporting seven processes—four articular, two transverse, and one spinous. There are exceptions to these general characteristics of a vertebra. For example, the first cervical vertebra (atlas vertebra) has neither a body nor spinous process. In addition, the second cervical vertebra (axis vertebra) has an odontoid process, which is a strong, prominent process, shaped like a tooth, rising perpendicularly from the upper surface of the body of the axis vertebra. Further details regarding the construction of the spine may be found in such common references as Gray's Anatomy, Crown Publishers, Inc., 1977, pp. 33-54, which is herein incorporated by reference.\nThe human vertebrae and associated connective elements are subjected to a variety of diseases and conditions which cause pain and disability. Among these diseases and conditions are spondylosis, spondylolisthesis, vertebral instability, spinal stenosis and degenerated, herniated, or degenerated and herniated intervertebral discs. Additionally, the vertebrae and associated connective elements are subject to injuries, including fractures and torn ligaments and surgical manipulations, including laminectomies.\nThe pain and disability related to the diseases and conditions often result from the displacement of all or part of a vertebra from the remainder of the vertebral column. Over the past two decades, a variety of methods have been developed to restore the displaced vertebra to their normal position and to fix them within the vertebral column. Spinal fusion is one such method. In spinal fusion, one or more of the vertebra of the spine are united together (“fused”) so that motion no longer occurs between them. The vertebra may be united with various types of fixation systems. These fixation systems may include a variety of longitudinal elements such as rods or plates that span two or more vertebrae and are affixed to the vertebrae by various fixation elements such as wires, staples, and screws (often inserted through the pedicles of the vertebrae). These systems may be affixed to either the posterior or the anterior side of the spine. In other applications, one or more bone screws may be inserted through adjacent vertebrae to provide stabilization.\nAlthough spinal fusion is a highly documented and proven form of treatment in many patients, there is currently a great interest in surgical techniques that provide stabilization of the spine while allowing for some degree of movement. In this manner, the natural motion of the spine can be preserved, especially for those patients with mild or moderate disc conditions. In certain types of these techniques, flexible materials are used as fixation rods to stabilize the spine while permitting a limited degree of movement."}
{"text": "1. Field of the Invention\nThe present invention relates to a printer for receiving print data transmitted from a client PC via a network and for effecting print processing on the basis of the received print data, as well as a method of controlling the same.\n2. Description of the Related Art\nIn printing making use of a network, since the client side and the printer are located away from each other, it is not easily possible to directly ascertain the state of the printer. When a print request has been issued from the client side to the printer, in a case where printing could not be effected properly due to a failure occurring in the printer, running out of paper, running out of toner, or the like, a user is unable to notice that printing could not be done. In addition, if the printer does not have a spool (HDD, a memory, etc.), the sent print data is inevitably deleted.\nAt the time of executing printing again, it has been necessary to prepare the print data again on the client side, and make a print request again to the printer.\nJapanese Patent Application Laid-open No. 2004-102916 discloses a technique in which electronic mail to which data subject to printing is attached and in which printing attributes, an address of a transfer destination, and the like are described in the text of the mail is sent from a personal computer to a printer, and in the event that an obstacle hampering printing has occurred in the printer, the received mail is transferred to the address of the transfer destination described in the text of the received mail, and a notification is given thereto to the effect that it has been transferred.\nHowever, particularly in the printer which does not have a spool, in the event that an obstacle hampering printing has occurred, the print data is inevitably deleted in the printer. At the time of executing printing again, it has been necessary to prepare the print data again on the client side, and subsequently issue a print request again.\nIn addition, even if the printer has a spool, in a case where the print data transmitted from the client side is large, the print data cannot be stored in the spool. Furthermore, if preceding print data is already stored in the spool, there are cases where following print data cannot be stored, in which case the print data is deleted.\nIn addition, according to the technique of Japanese Patent Application Laid-open No. 2004-102916, it is possible to know that printing could not be effected properly, but since only a notification is given to the effect that the print data has been transferred, in order for the client to obtain a printing result again, it is necessary to send from the client side electronic mail with the print data attached thereto. This operation is troublesome, and the client side is inevitably restrained until printing is completed.\nAccordingly, an object of the present invention is to provide a printer which, in the event of printing failure by the printer, makes it possible for the client device to obtain the print data without forming the print data again, as well as a method of controlling the same."}
{"text": "Fabric items such as clothing, towels, bedding, etc. can be colored using a variety of different dyes and dyeing processes. In a residential setting, caring for these dyed fabric items may present consumers with several challenges. Some dyed fabric items may have excess or loose dye that can wash off during a normal wash cycle in a clothes washer and redeposit on other items in the laundry load or bleed onto differently dyed areas of the same item, for example. Excess or loose dyes may also rub off onto the consumer or other surfaces during wear or use. Sorting the laundry items before washing into loads of “like color” or washing items separately may address some dye transfer concerns, but can by time consuming and inefficient for the user. In addition, mistakes in sorting loads can lead to dye transfer which cannot be easily removed, potentially ruining the item."}
{"text": "Vehicles may be equipped with bumper systems and crash protection structures that plastically deform to absorb energy in the event of an impact. When an object impacts a vehicle at a position that is offset from the centerline of the vehicle such that the object overlaps a portion of the bumper, the ability of all of the energy absorbing structures of the vehicle to absorb energy associated with the impact may be reduced. In some impact configurations, the energy absorbing structures of the vehicle may not be activated or may be only partially activated because the object does not come into contact or only partially comes into contact with associated bumper or vehicle structures. Therefore, the bumper and the energy absorbing structures of the vehicle may have a reduced effect on the dissipation of the energy of the impact. Instead, the energy from the impact may be directed into various vehicle structures.\nIn one example, a substantial portion of energy from an impact with a small front bumper overlap may be directed into an upper side member that is positioned forward of the vehicle A-pillar. As energy is directed into the upper side member, the upper side member may rotate inboard, deflecting away from the location of the impact. When the upper side member deflects away from the location of the impact, the upper side member may absorb less energy of the impact than when the upper side member does not deflect away from the location of the impact.\nAccordingly, a need exists for alternative structures for transferring energy and absorbing energy from a small front bumper overlap impact."}
{"text": "1. Field\nThe present invention generally relates to induction RF fluorescent light bulbs, and more specifically to reduction of electromagnetic interference from an induction RF fluorescent light bulb with a ferromagnetic core.\n2. Description of Related Art\nDischarge lamps create light by exciting an electrical discharge in a gas and using that discharge to create visible light in various ways. In the case of fluorescent lamps the gas is typically a mixture of argon, krypton and/or neon, plus a small amount of mercury. Other types of discharge lamps may use other gasses. The gas is contained in a partially evacuated envelope, typically transparent or translucent, typically called a bulb or arc tube depending upon the type of lamp.\nIn conventional discharge lamps electrically conductive electrodes mounted inside the bulb or arc tube along with the gas provide the electric field used to drive the discharge.\nUse of electrodes can create certain problems. First, the discharge is typically designed to have a relatively high voltage in order to minimize losses at the electrodes. In the case of fluorescent lamps, this may lead to long, thin lamp structures, which function well for lighting office ceilings, but are not always a good fit for replacing conventional incandescent lamps. Fluorescent lamps designed to replace incandescent lamps, known as compact fluorescent lamps, or CFLs, are typically constructed by bending the long, thin tube, such as into multiple parallel tubes or into a spiral, which is now the most common form of CFLs. A plastic cover shaped like a conventional incandescent lamp is sometimes placed over the bent tubes to provide a more attractive shape, but these covers absorb light, making the lamp less efficient. Bent and spiral tube lamps also have wasted space between the tubes, making them larger than necessary. The use of a cover increases the size further.\nThe use of electrodes can create problems other than shape and size. Electrodes can wear out if the lamp is turned on and off many times, as is typical in a residential bathroom and many other applications. The life of the electrodes can also be reduced if the lamp is dimmed, because the electrodes are preferably operated in a specific temperature range and operation at different power levels can cause operation outside the preferred ranges, such as when operating at lower power, which can allow the electrodes to cool below the specified temperature range.\nIn addition, the long thin shape selected, because it is adapted to allow use of electrodes, tends to require time for mercury vapor to diffuse from one part of the tube to another, leading to the long warm-up times typically associated with many compact fluorescent lamps.\nFinally, the electrodes are normally designed to be chemically compatible with the gas used in the lamp. While this is not usually a concern with typical fluorescent lamps, it can be a problem with other types of discharge lamps.\nOne way to avoid the problems caused by electrodes is to make a lamp that does not use electrodes, a so-called electrodeless lamp. In an electrodeless lamp, the discharge may be driven by, for example, 1) an electric field created by electrodes mounted outside the bulb or arc tube; 2) an electric field created by a very high frequency electromagnetic field, usually in combination with a resonant cavity, or 3) an electric field created by a high frequency magnetic field without the use of a resonant cavity. This latter lamp is called an induction-coupled electrodeless lamp, or just “induction lamp.”\nIn an induction lamp, a high frequency magnetic field is typically used to create the electric field in the lamp, eliminating the need for electrodes. This electric field then powers the discharge.\nSince induction lamps do not require use of electrodes, they do not need to be built into long thin tubes. In fact, a ball-shaped bulb, such as the bulb used for conventional incandescent lamps, is a preferred shape for an induction lamp. In addition, since induction lamps do not use electrodes, they can be turned on and off frequently without substantial adverse impact on loss of life. The absence of electrodes also means that induction lamps can be dimmed without reducing lamp life. Finally, the ball-shaped lamp envelope allows rapid diffusion of mercury vapor from one part of the lamp to another. This means that the warm-up time of induction lamps is typically much faster than the warm-up time of most conventional compact fluorescent lamps.\nInduction lamps fall into two general categories, those that use a “closed” magnetic core, usually in the shape of a torus, and those that use an “open” magnetic core, usually in the shape of a rod. Air core induction lamps fall into this latter category. Closed core lamps are usually operated at frequencies generally above 50 kHz, while open core lamps usually require operating frequencies of 1 MHz and above for efficient operation. The lower operating frequency of closed core induction lamps makes them attractive; however, the bulb design required to accommodate the closed core makes them generally unsuitable for replacing standard in incandescent lamps. Open core induction lamps, while requiring higher operating frequencies, allow the design of lamps that have the same shape and size as common household incandescent lamps. This disclosure is primarily is addressed to the open core category of induction lamps.\nIn spite of their obvious advantages, there are very few open core induction lamps on the market today. One reason for the lack of commercially successful products is the cost of the high frequency ballast. Conventional fluorescent lamps, including CFLs, can be operated at frequencies from 25 kHz to 100 kHz, a frequency range where low cost ballast technology was developed in the 1990s, and closed core induction lamps can be operated at frequencies from 50 kHz to 250 kHz, for which the ballasts are only slightly more expensive. However, open core induction lamps typically require operating frequencies of 1 MHz or higher. The United States Federal Communications Commission (FCC) has established a “lamp band” between 2.51 MHz and 3.0 MHz that has relaxed limits on the emission of radio frequency energy that may interfere with radio communication services. Cost effective open core induction lamps may preferably have an operating frequency of at least 2.51 MHz.\nThe lack of commercially successful open core induction lamps may be traced to the absence of a low cost ballast that can operate in the 2.51 MHz to 3.0 MHz band while meeting all the requirements of the FCC; that is small enough to fit into a lamp; that has a ballast housing that is the same size and shape as a conventional incandescent lamp; and that can be dimmed on conventional TRIAC dimmers found in homes in certain major markets, such as the U.S. The present disclosure addresses one or more of these issues. Therefore a need exists for improved induction lamps, especially in residential applications."}
{"text": "Many retail stores, whether online and/or brick and mortar, participate in large and expensive marketing plans in order to develop, to grow, and to maintain a customer base. Part of the marketing plan may include an advertising plan that may include offers to new and/or existing customers.\nThe offers may include coupons (for example, a coupon for a discount), rebates, product recommendations, indication of a price drop, and the like. The offers may be time sensitive, e.g., valid for a limited time or during a predetermined time period. The offers may be associated with goods or services that are limited in nature. For example, an offer may be associated for a product that the retail store only has in its possession a limited quantity. Or, the retail store may only be financially able to provide the goods or services under the offer on a limited basis.\nFurthermore, it may be cost prohibitive to provide the offers to a large number (or all) of a set of customers. One downside to providing the offer to a large number of customers is that such a widespread offer may be expensive to provide. Also, the likelihood of a large number of the customers that utilize the offer may be quite small.\nThe present invention is aimed at one or more of the problems identified above."}
{"text": "1. Field of the Invention\nThe present invention relates to a medical instrument such, for example, as a forceps, a scissor-type instrument or needle holder, with a pair of distal jaws, a proximal handle which actuates gripping motion to the jaws, and a means for limiting the holding force of the jaws to a predetermined maximum value.\n2. Description of the Prior Art\nGerman Patent document DE-A 36 01 166 discloses a forceps for gripping tissue of a patient. The forceps has a fixed jaw and a movable jaw, the fixed jaw being connected to a fixed grip and the movable jaw being connected to a pivotable grip. To limit the compressive force that the jaws can exert on the tissue, the pivotable grip is divided into two pieces --an upper and a lower piece. The two pieces are articulated or joined to one another and supported against each other by a spring member. In operation, when the applied pressure exceeds a predetermined value, the lower piece becomes deflected relative to the upper piece so that the pivotable grip can no longer pivot. However, upon further movement of the pieces towards each other, the holding force of the jaws increases, depending on the stiffness of the spring member. Consequently, this increased holding force could ultimately damage, perforate or sever the tissue held between the jaws.\nGerman Patent document DE-A 40 10 775 discloses a surgical forceps having a fixed jaw and a movable jaw, which are connected to a fixed grip and a pivotable grip, respectively. The pivotable grip consists of two parts in alignment with each other. The two parts are unreleasably connected to each other by a spiral spring element. Detent elements or pawls are provided on the two parts of the pivotable grip, which are releasably engaged with each other. The holding force applied by the pivotable grip cannot be limited to a predetermined value. Consequently, the surgical forceps can damage blood vessels, tissue, or organs if a user inadvertently applies an excessive force during a surgical operation or procedure."}
{"text": "Mechanical and chemical-mechanical planarizing processes (collectively \"CMP\") are used in the manufacturing process of microelectronic devices to form a flat surface on semiconductor wafers, field emission displays, and many other microelectronic-device substrates and substrate assemblies. FIG. 1 is a partially schematic, isometric view of a conventional web-format planarizing machine 10 that has a platen 20. A sub-pad 11 is attached to the platen 20 to provide a flat, solid workstation for supporting a portion of a web-format polishing pad 16 in a planarizing zone \"A\" during planarization. The polishing pad 16 has a rear surface 19 that engages the sub-pad 11 and a planarizing surface 18 facing opposite the rear surface 19 to planarize a substrate 12.\nThe planarizing machine 10 also has a pad-advancing mechanism, including a plurality of rollers, to guide, position and hold the polishing pad 16 over the sub-pad 11. The pad-advancing mechanism generally includes a supply roller 24, first and second idler rollers 21a and 21b, first and second guide rollers 22a and 22b, and a take-up roller 23. As explained below, a motor (not shown) drives the take-up roller 23 and the supply roller 24 to advance and retract the polishing pad 16 over the sub-pad 11 along a travel path T--T. The first idler roller 21a and the first guide roller 22a press an operative portion of the polishing pad 16 against the sub-pad 11 to hold the polishing pad 16 stationary during operation.\nThe planarizing machine 10 further includes a carrier assembly 30 to translate the substrate 12 over the polishing pad 16. In one embodiment, the carrier assembly 30 has a head 31 to pick up, hold and release the substrate 12 at appropriate stages of the planarizing process. The carrier assembly 30 also has a support gantry 32 and a drive assembly 33 that can move along the gantry 32. The drive assembly 33 has an actuator 34, a drive shaft 35 coupled to the actuator 34, and an arm 36 projecting from the drive shaft 35. The arm 36 carries the head 31 via a terminal shaft 37. The actuator 34 orbits the head 31 about an axis B--B (as indicated by arrow R.sub.1) and can rotate the head 31 about an axis C--C (as indicated by arrow R.sub.2) to move the substrate 12 over the polishing pad 16 while a planarizing fluid 17 flows from a plurality of nozzles 38 in the head 31. The planarizing fluid 17 may be a conventional CMP slurry with abrasive particles and chemicals that etch and/or oxidize the surface of the substrate 12, or the planarizing fluid 17 may be a non-abrasive planarizing solution without abrasive particles. In most CMP applications, conventional CMP slurries are used on conventional polishing pads, and planarizing solutions without abrasive particles are used on fixed-abrasive polishing pads.\nIn the operation of the planarizing machine 10, the carrier assembly 30 presses the substrate 12 against the planarizing surface 18 of the polishing pad 16 as the carrier head 31 moves the substrate 12 over the planarizing surface 18. The polishing pad 16 moves across the sub-pad 11 along the pad travel path T--T either during or between planarizing cycles to change the particular portion of the polishing pad 16 in the planarizing zone A. For example, the supply and take-up rollers 24, 23 can drive the polishing pad 16 between planarizing cycles such that a point P moves incrementally across the sub-pad 11 to a number of intermediate locations I.sub.1, I.sub.2, etc. Alternatively, the rollers 24, 23 may drive the polishing pad 16 between planarizing cycles such that the point P moves all the way across the sub-pad 11 toward the take-up roller 23 to completely remove a used or post-operative portion of the polishing pad 16 from the planarizing zone A. The rollers 24, 23 may also continuously drive the polishing pad 16 at a slow rate during a planarizing cycle such that the point P moves continuously across the sub-pad 11 during planarization.\nThe planarizing machine 10 can also include a planarizing surface cleaner 40 (shown schematically in FIG. 1) positioned between the platen 20 and the take-up roller 23 to clean the post-operative portion of the polishing pad 16. The planarizing surface cleaner 40 can include a brush 41 having bristles that contact the planarizing surface 18 of the polishing pad 16 and a liquid dispenser 42 positioned proximate to the brush 41 to dispense a cleaning liquid on the planarizing surface 18. Accordingly, the planarizing surface cleaner 40 can clean the post-operative portion of the polishing pad 16 as it moves off the platen 20 along the travel path T--T. Once the post-operative portion of the polishing pad 16 has been cleaned, it can be translated back onto the platen 20 along the travel path T--T and into the planarizing zone A for another planarizing cycle.\nOne drawback with the apparatus 10 shown in FIG. 1 is that the rear surface 19 of the polishing pad 16 can become contaminated with debris (such as liquid and/or particulate matter) during the planarizing process and/or the cleaning process. The debris can become trapped between the polishing pad 16 and the sub-pad 11, causing a local bump or other non-uniformity to form in the planarizing surface 18. The non-uniformity in the planarizing surface 18 can create a non-uniformity in the substrate 12 and/or can cause the polishing pad 16 to wear in a non-uniform manner.\nA further drawback is that liquid on the rear surface 19 of the polishing pad 16 can form an adhesive bond between the polishing pad 16 and the sub-pad 11. The adhesive bond can inhibit relative movement between the polishing pad 16 and the sub-pad 11 when the polishing pad 16 moves along the travel path T--T. In one conventional method, the idler rollers 21a, 21b and/or the guide roller 22a move the polishing pad 16 normal to the upper surface of the sub-pad 11 to break the adhesive bond. However, the action of the rollers against the polishing pad 16 may not be effective to separate the polishing pad 16 from the sub-pad 11. Furthermore, if the polishing pad 16 is dragged over the sub-pad 11, the frictional contact between the two can abrade particulate matter from the polishing pad 16 and/or the sub-pad 11, which can cause a bump or other non-uniformity to form in the planarizing surface 18, as discussed above."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nFirearms conventionally incorporate a sight to aid in aligning a trajectory of the firearm with a target. In one configuration, the sight includes an upwardly extending arm fixed relative to a barrel of a firearm, whereby a user of the firearm may properly align an end of the barrel with a target by aligning the upwardly extending arm with the target.\nIn addition to use of a fixed, upwardly extending arm, conventional firearm sights may also incorporate an optical element that displays an illuminated reticle for use in aligning a barrel of a firearm with a target. One such prior-art sight is disclosed in U.S. Pat. No. 6,327,806. The foregoing sight incorporates an optical element, which receives light from a light emitting diode (LED) and displays a reticle on a lens for use by a user in aligning a barrel of a firearm with a target. Such a sight incorporating a lens and an illuminated reticle is generally an improvement over a firearm incorporating a fixed, upwardly extending arm, as the illuminated reticle of the optical sight may be viewed from numerous angles from a rear portion of the firearm and does not have to be exactly aligned with an eye of the user. Allowing the reticle to be viewed from numerous angles from an area generally behind the firearm allows the user to be positioned somewhat offset from a longitudinal axis of the firearm while still maintaining a barrel of the firearm trained on a target.\nWhile the foregoing optical sight is an improvement over a fixed, upwardly extending arm disposed proximate to an end of a firearm, conventional optical sights typically include an optical lens having a generally convex upper surface, which is easily fractured if dropped. While conventional optical sights typically include a housing having a portion extending over the convex upper surface, the housing typically includes a similar convex shape and, as such, transmits a force applied at an outer surface thereof directly to an outer surface of the lens, thereby causing the lens to fracture. Once the lens of the optical sight is fractured, the sight may not be used and, therefore, reduces the overall effectiveness of the firearm.\nIn addition to the likelihood of fracture, conventional optical sights suffer from the disadvantage of including an LED, which requires a power source to illuminate a reticle. While such LEDs adequately illuminate a reticle, the power source supplying power to the LED is not infinite. Because the power source supplying power to the LED is not infinite, care must be taken to routinely check the life of the power source to ensure that the reticle is consistently displayed. While recharging or replacing the power source of a conventional sight is relatively simple, such tasks become difficult in a military or law enforcement operation where time is of the essence and reliability on equipment is key. Having a power source expire during a law enforcement or military operation reduces the overall effectiveness of the firearm on which the optical sight is mounted and, as a result, reduces the effectiveness of the law enforcement agent or soldier.\nBased on the foregoing, an optical sight incorporating multiple light sources, such as, for example, an LED, a fiber optic, and a tritium lamp, that can accommodate various ambient-light conditions is desirable in the industry. Incorporating multiple light sources into an optical sight provides flexibility in illuminating a reticle, as each source or a combination of sources can be chosen based on the particular ambient-light conditions.\nIn one configuration, light from the fiber optic and tritium lamp may be combined to illuminate a reticle. In another configuration, light from the LED may additionally or alternatively be used should the supplied light from the fiber optic and/or tritium lamp be insufficient. Further yet, light from any one of the sources may be used independently of the other sources. In any of the foregoing configurations, providing an optical sight with multiple light sources allows the optical sight to be used in virtually any ambient-light condition and provides the user with a reliable and useful sight."}
{"text": "The present invention relates to photocopying systems and, more particularly, to developing apparatus for developing light sensitive paper and polyester film. The development process requires the sensitive side of the copy material to be coated with a very small amount of developer liquid, this liquid to be distributed evenly over the surface of the copy material, the quantities being in the order of 1.75 to 2.25 grams of liquid per square meter of paper. The major advantage of applying the liquid in such small quantities is that the copies emerge direct from the developing apparatus dry to the touch.\nOne existing system known as the \"pressure diazo\" system is disclosed in British Pat. Nos. 1,215,705, 1,349,597, 1,349,598, 1,349,599 and U.S. Pat. Nos. 3,626,833 and 3,640,203. Basically the \"pressure diazo\" system consists of a rotating roller with a rubber surface having a fine surface texture. The rubber surface effectively transports developer liquid from the developer supply to the surface of the copy material as they meet in rolling contact with each other. The quantity of developer transported is limited by a metering blade which effectively only allows the small amounts of liquid held within the \"valleys\" of the surface texture to reach the copy material. The transfer of developer liquid from the rubber surface to the copy material is assisted by a pressure applicator device which serves to press the copy material firmly against the rotating rubber surface.\nBoth the metering device and the pressure applicator device are blade members and are mounted in such a way that, under the influence of a control mechanism, they can be moved into and out of engagement with the roller surface. When engaged the metering blade traps a reservoir of developer liquid downstream of the pressure applicator blade. When a print is presented to the unit, the copy material and the developer laden roller surface meet and are pressed together under the influence of the pressure applicator blade, this causes the major part of the liquid in the \"valleys\" to transfer onto the activated side of the copy material thus causing it to develop.\nWhen the machine is in operation and the roller is rotating, but the unit is not processing copy material, the pressure of both blades acts directly onto the roller surface. The points of contact are lubricated to some extent by the developer liquid, but this degree of lubrication is not adequate when the roller is rotating at the minimum speed. In order to reduce contact friction and overcome this problem, the blades have been coated with a fluoro carbon plastic such as tetrafluoroethylene polymer.\nThe movement of the blades, into and out of engagement with the roller surface, occurs when the machine is switched on or off respectively, which ensures that the blades do not bear down under pressure onto the surface of the roller when it is stationary. The effect of this if it should occur, for more than a few hours, would be to mark the surface of the roller with a smooth line measuring the same length as that of the blade. On a microscope scale, the \"crests\" of the surface texture become depressed into the space of the \"valleys\" of the same surface texture, thus reducing the liquid carrying capacity of the surface at the pressure set point. When the machine is then used normally, a line of poor development on the copy material will appear corresponding to the line on the roller.\nIt is therefore the principal object of the present invention to provide a developing apparatus in a photocopying system which overcomes the disadvantages of the prior art by including means to apply pressure on the copy paper against the surface of the application roller without causing marks on the surface of the roller and without requiring means to move the pressure application means away from the roller when the roller is stationary.\nOther objects, features and advantages of the invention will become apparent from the description of the invention in connection with the accompanying drawing to be described more fully hereinafter."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid lens containing liquids for changing focus.\n2. Description of the Related Art\nIn recent years, remarkably high performance owing to the increased number of pixels and the like in a camera mounted to a portable terminal device or the like has been attained, and there has been a demand for a highly accurate optical zooming function along with the high performance of the camera. In ordinary optical zooming, it is necessary to change a focal length between an image sensor (e.g., a charge-coupled device (CCD) sensor) in a camera and a lens, and the focal length is changed by mechanically moving the lens. Therefore, in the case of mounting the ordinary optical zooming, a space for the movement of the lens is required. However, there are demands for downsizing and thinning of a camera in portable terminal devices, and an increase in size due to the space for movement of a lens should be avoided even in the case of mounting the optical zooming function.\nAs a method for attaining the space saving, Japanese Patent Application Laid-Open No. 01-302301 discusses a liquid lens that does not require the space for movement of the lens.\nIn Japanese Patent Application Laid-Open No. 01-302301, the lens includes two phase liquids which are partitioned by an elastic transparent film, and the shape of the elastic transparent film is changed by adjusting an amount of each of the two phase liquids by an externally connected pump, thereby adjusting a lens focus. In the liquid lens, since the lens focus is changed when the shape of the liquid lens itself is changed by the liquid amount adjustment, the space for movement of the lens is no longer necessary.\nAlso, as another mode of the liquid lens, Japanese Patent Application Laid-Open No. 2000-249813 discusses a lens including a single phase liquid which is covered with an elastic transparent film, of which a lens focus is changed by changing the shape of the elastic film by an actuator.\nHowever, though the space for movement of the lens is not required in the liquid lens discussed in Japanese Patent Application Laid-Open No. 01-302301, the pump for adjusting the liquid amounts and a liquid reserve tank are required. Therefore, it is difficult to attain the space saving of an entire device to which the liquid lens is to be mounted. Also, though the pump and the tank are not required in the liquid lens discussed in Japanese Patent Application Laid-Open No. 2000-249813, it is in some cases impossible to ensure satisfactory accuracy of the focus adjustment because of deflection of the elastic film due to an external cause, such as gravity or external impact, since the elastic film is exposed to the outside."}
{"text": "In a typical residential environment, the hot water heater is responsible for a large portion of a household's power consumption bill. Therefore, for a utility cost management plan to succeed, the hot water heater must be considered.\nIn order to save on the hot water bill, the hot water heater should be in top condition. A hot water heater that is electrically operated, typically has a one or two heating elements protruding into the water tank. As electricity is introduced, the elements become hot and heat the water within the tank. In order to efficiently operate the water heater, these elements must be kept in good shape.\nHowever much the public desires to believe to the contrary, the residential water supply is not pure water. The water will typically contain impurities such as minerals, sediments, sand, debris, and the like. These impurities are especially pronounced in households that rely on well water and those located in coastal areas.\nThese impurities will settle on the bottom of the hot water tank and accumulate over time. The accumulation of sediment builds up on and cakes to the elements, especially the element located close to the bottom of the tank. A sediment buildup on the element decreases its efficiency and increases the cost of heating the water within the tank. As the accumulation of the sediment continues and the buildup on the elements increases, the elements become less and less efficient until they fail, necessitating a replacement. The sediment accumulation also decreases the useful life of the hot water tank itself.\nWhat is needed is a way to be able to remove the sediment without adding to the purchase cost of the hot water tank assembly. One simple solution is to periodically remove the sediment, thereby preventing an accumulation within the tank and a buildup on the elements. However, the only opening available into present day hot water tanks is through the opening through which the element is inserted, a very small opening indeed. With such a small opening, a person desiring to remove debris would be required to use a small instrument such as a teaspoon. Such a process would be slow and inefficient and would discourage all but the extremely dedicated.\nWhat is needed is a hot water tank assembly that has an opening in the tank that will permit a person to insert his hand directly into the tank. Such an assembly would give a person easy access into the interior of the hot water tank and would permit that person to remove accumulated sediment quickly and easily. With regular sediment removal, the elements would not experience a sediment buildup and would operate more efficiently as well as experience a longer life span. The result would a lower hot water bill and lower maintenance costs of both the elements and the tank assembly itself."}
{"text": "Over the last decade, there has been a substantial increase in the use and deployment of network enabled client devices. These client devices may join a Layer 2 (L2) local network in which the client devices operate on the same subnet (i.e., address range). Since client devices operate using the same subnet, communications between client devices may be performed using media access control (MAC) addresses. In particular, network devices within this L2 local network may use Address Resolution Protocol (ARP) requests and responses to resolve a mapping between Internet Protocol (IP) addresses and MAC addresses for client devices. Based on this mapping, client devices may communicate with each other within the local network using MAC addresses.\nIn some situations, the local network described above may be communicatively coupled to one or more geographically remote networks such that client devices in the local network may communicate with a public cloud services to obtain access to particular resources, such as applications and/or storage. Examples of public clouds include Amazon® AWS, Google® Compute Engine and Microsoft® Azure Services Platform. In this configuration, the local network typically operates on a different subnet (i.e., address range) than the remote networks, producing a resultant Layer 3 (L3) network. As a L3 network, client devices in the local network communicate using intermediate routing devices instead of directly using MAC addresses as with a L2 network. Accordingly, the client devices on the local network and the remote networks are in different broadcast domains and may not efficiently transmit multicast or broadcast packets in this L3 network. However, as described in U.S. Pat. No. 9,825,906, a mechanism has been developed to manage, create and connect a local network to private cloud networks within a public cloud to collectively operate as a L2 network. However, this deployment may experience throughput performance limitations, given limited processing, memory and network bandwidth resources. A solution to solve the performance limitations is needed."}
{"text": "Certain vehicles, such as aircraft, can be operated in an autopilot mode, in which an autopilot system controls various systems of the aircraft to control the path of the aircraft under the supervision of a pilot. In certain instances, the autopilot system may be used to direct the vertical trajectory of the aircraft, such as a vertical descent of the aircraft for landing. Generally, the vertical trajectory of an aircraft is defined as a function of altitude with associated speed constraints along the aircraft lateral distance to the end point of the leg. While controlling the vertical descent of the aircraft, the autopilot system may encounter conditions where the aircraft may deviate from the planned vertical path due to varying wind conditions or incorrect energy management in descent or approach. Given the pilot's generally high workload during descent and approach for landing, the pilot may not be able to easily discern the new path of the aircraft taken by the autopilot system.\nAccordingly, it is desirable to provide improved systems and methods for providing situational awareness of the current and the future state of a vehicle, such as an aircraft, during an autopilot mode, by displaying the current and the future state on a display associated with the vehicle to enable the pilot to discern the current and future path of the vehicle. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "1. Field of the Invention\nThe present invention relates to a fixing device and a fixing method for fixing a toner image on a recording sheet.\n2. Description of the Related Art\nIn Japanese Patent Application Publication No. 2005-274948, a fixing device including a heat roller having a heat generation surface for generating heat for heating a recording sheet by heat generated from the heat generation surface, a pressure roller for pressing the recording sheet to the heat generation surface of the heat roller, and a thermostat for stopping heat generation of the heat generation surface when the temperature at the monitoring portion where the peak of the temperature appears exceeds a predetermined threshold value in the second area of the heat generation surface of the heat roller excluding the first area where the recording sheet always makes contact with at time of press-fit by the pressure roller is disclosed.\nFurther, in Japanese Patent Application Publication No. 2003-43854, an image forming apparatus composed of a fixing device including a first heater arranged in the neighborhood of the central part of fixing rollers in the axial direction, a second heater arranged in the neighborhood of both ends thereof, a first temperature detecting means for detecting the temperature of the first heater having larger power than that of the second heater, and a second temperature detecting means for detecting the temperature of the second heater is disclosed.\nTo heat the central part and end parts of each recording sheet, separate heaters are used and to detect abnormal temperatures of the respective heaters, at least two thermostats are required. However, due to use of two thermostats, a problem arises that the whole cost of the copier is increased and no inexpensive products can be provided to users."}
{"text": "Lactobacillus acidophilus is a Gram-positive, rod-shaped, non-spore forming, homofermentative bacterium that is a normal inhabitant of the gastrointestinal and genitourinary tracts. Since its original isolation by Moro (1900) from infant feces, the “acid loving” organism has been found in the intestinal tract of humans, breast-fed infants, and persons consuming high milk, lactose, or dextrin diets. Historically, Lactobacillus acidophilus is the Lactobacillus species most often implicated as an intestinal probiotic capable of eliciting beneficial effects on the microflora of the gastrointestinal tract (Klaenhammer and Russell (2000) “Species of the Lactobacillus acidophilus complex,” Encyclopedia of Food Microbiology, 2:1151-1157. Robinson et al., eds. (Academic Press, San Diego, Calif.). Lactobacillus acidophilus can ferment hexoses, including lactose and more complex oligosaccharides, to produce lactic acid and lower the pH of the environment where the organism is cultured. Acidified environments (e.g., food, vagina, and regions within the gastrointestinal tract) can interfere with the growth of undesirable bacteria, pathogens, and yeasts. The organism is well known for its acid tolerance, survival in cultured dairy products, and viability during passage through the stomach and gastrointestinal tract. Lactobacilli and other commensal bacteria, some of which are considered probiotic bacteria that “favor life,” have been studied extensively for their effects on human health, particularly in the prevention or treatment of enteric infections, diarrheal disease, prevention of cancer, and stimulation of the immune system. Lactobacilli have also been studied for their influence on dairy product flavor, and functional and textural characteristics. Genetic characterization of other Lactobacillus species (e.g., L. johnsonii and L. rhamnosus) has been described (see e.g., U.S. Pat. No. 6,476,209; U.S. Pat. No. 6,544,772; U.S. Patent Publication Nos. 20020159976, 2003013882 & 20040009490; PCT Publication No. WO 2004/031389; PCT Publication No. 2003/084989; PCT Publication No. WO 2004/020467).\nBacterial growth requires specific transport systems to import nutrients from the external environment. Lactic acid bacteria transport molecules into and out of the cell via three systems: primary transport, secondary transport, and group translocation. In primary transport, chemical (primarily ATP), electrical, or solar energy is used to drive transport. ATP-binding cassette (ABC) transporters are the most abundant class of primary transport systems in lactic acid bacteria. In this system, ATP hydrolysis is linked with substrate translocation across the membrane for both the import of sugars and compatible solutes and the export of products such as drugs or toxins that are undesirable to the cell, or cellular components that function outside of the cell, such as cell wall polysaccharides. In general, ABC transporters are relatively specific for their substrates, but some are multispecific, such as the multidrug transporters.\nSecondary transport systems use electrochemical gradients to provide the energy for sugar translocation. They comprise symporters, which cotransport two or more solutes, uniporters, which transport one molecule, and antiporters, which countertransport two or more solutes. Symporters generally couple the uphill movement of the substrate to the downhill movement of a proton (or ion), antiporters use the ion gradient for excretion of a product, and uniporters do not use a coupling ion (Poolman (2002) Antonie van Leeuwenhoek 82:147-164).\nGroup translocation involves the phosphoenolpyruvate (PEP)-dependent phosphotransferase system (PTS), which couples the uptake of a carbohydrate or alditol with its phosphorylation (Poolman (2002), supra). The phosphate group originates from the conversion of PEP into pyruvate, and the subsequent phosphorylation involves the energy coupling proteins, Enzyme I and HPr, as well as substrate-specific phosphoryl transfer proteins IIA, IIB and IIC.\nMultidrug transporters may be separated into two major classes, secondary multidrug transporters and ABC transporters. Secondary multidrug transporters may be further divided into distinct families, including the major facilitator superfamily (MFS), the small multidrug resistance family (SMR), the resistance-nodulation-cell division family (RND), and the multidrug and toxic compound extrusion family (MATE) (Putman et al. (2000) Microbiol. Mol. Biol. Reviews 64:672-693). Secondary multidrug transporters use the electrochemical gradients, as described herein, to extrude drugs from the cell. ABC-type multidrug transporters use energy from ATP hydrolysis to pump drugs out of the cell (Putman et al. (2000), supra).\nBacteria are able to metabolize various carbohydrates by utilizing transport proteins and enzymes with different carbohydrate specificities, in addition to employing diverse regulatory mechanisms, such as catabolite repression. The isolation and characterization of these proteins allows for the development of essential probiotic products with numerous applications, including those that benefit human and/or animal health, and those concerned with food production and safety. The proteins can also be used in developing transgenic plants with altered growth or survival capabilities."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nLayered heaters are typically used in applications where space is limited, when heat output needs vary across a surface, where rapid thermal response is desirous, or in ultra-clean applications where moisture or other contaminants can migrate into conventional heaters. A layered heater generally comprises layers of different materials, namely, a dielectric and a resistive material, which are applied to a substrate. The dielectric material is applied first to the substrate and provides electrical isolation between the substrate and the electrically-live resistive material and also reduces current leakage to ground during operation. The resistive material is applied to the dielectric material in a predetermined pattern and provides a resistive heater circuit. The layered heater also includes leads that connect the resistive heater circuit to an electrical power source, which is typically cycled by a temperature controller. The lead-to-resistive circuit interface is also typically protected both mechanically and electrically from extraneous contact by providing strain relief and electrical isolation through a protective layer. Accordingly, layered heaters are highly customizable for a variety of heating applications.\nLayered heaters may be “thick” film, “thin” film, or “thermally sprayed,” among others, wherein the primary difference between these types of layered heaters is the method in which the layers are formed. For example, the layers for thick film heaters are typically formed using processes such as screen printing, decal application, or film dispensing heads, among others. The layers for thin film heaters are typically formed using deposition processes such as ion plating, sputtering, chemical vapor deposition (CVD), and physical vapor deposition (PVD), among others. Yet another series of processes distinct from thin and thick film techniques are those known as thermal spraying processes, which may include by way of example flame spraying, plasma spraying, wire arc spraying, and HVOF (High Velocity Oxygen Fuel), among others.\nIn some electrical heater applications, it may be desirable to vary the watt density of the heater in certain areas in order to tailor the amount of heat delivered to the specific part or device being heated or to account for inherent variations in heat distribution along the heater trace or element. Known electrical heaters typically vary the spacing of the resistive circuit pattern such that where the spacing is smaller and the trace of the resistive circuit pattern is closer, the watt density is higher, for a series circuit configuration. Conversely, the larger the spacing between the traces of the resistive circuit pattern, the lower the watt density in those regions. In other known electrical heaters, the width of the trace of the resistive circuit pattern is varied along its length in order to vary the watt density, wherein the wider the trace the lower the watt density and the narrower the trace the higher the watt density, again, for a series circuit configuration."}
{"text": "The present invention relates to an outboard jet drive boat and especially to an outboard jet drive having an engine and jet drive mounted in a housing which is removably attached to the transom of a boat hull.\nThere have been several proposed types of outboard set drives for watercraft but most are similar to an outboard motor in which the outboard motor propeller and lower unit have been replaced with a jet drive. The jet drive includes a jet pump in the lower unit that operates so as to provide propulsion force for a watercraft. There are advantages in employing jet pumps for propulsion units as opposed to propellers. The jet drive permits operation in shallower water and also the propeller is shrouded and there is less likelihood of injury. There has been a variety of proposed constructions for outboard jet drives for positioning the jet pump in different positions relative to the hull transom and bottom of the transom but in a typical jet drive, the engine and jet drive are located directly in the hull with an opening in the bottom of the hull for capturing water passing under the hull and then utilizing the jet pumps to thrust the water out the rear of the hull to propel the boat. Outboard jet drive units are made similar to typical outboard motors with a motor driving a drive unit which operates a jet drive unit.\nPrior art outboard liquid jet propulsion units can be seen in the Nanami U.S. Pat., No. 5,536,187, for an outboard jet drive for watercraft in which the jet propulsion unit is disposed forwardly of the transom and beneath the undersurface of the hull for improving its pumping efficiency while the motor is attached to the transom of the boat. In the Jordan U.S. Pat., No. 4,459,117, a liquid jet propulsion unit is driven by a conventional outboard motor. The drive of the motor directly rotates an impeller which draws water into the impeller chamber and through an outlet chamber and nozzle to drive the craft forward. In the Miyamoto U.S. Pat., No. 4,457,724, an apparatus for driving a surfboard includes an internal combustion engine mounted in a box which is mounted on the rear portion of the surfboard. A water jet propelling device is driven by the internal combustion engine for propelling the surfboard. The exhaust gas system in the water jet propelling device is positioned in the box. In the U.S. Pat. to Boyer et al., No. 4,942,838, an inflatable watercraft has a portable engine package. The engine package includes an internal combustion engine mounted in a thin fiberglass hull. The base plate of the hull includes a water inlet scoop for feeding water to the pump and an exhaust port for exhausting the water. The pumps high pressure water outlet is pointed in the aft direction above the water line to propel the craft by the reaction force resulting from the high velocity water jet. In the F.C. Clark U.S. Pat., No. 3,055,175, a marine propulsion unit takes a conventional outboard motor and replaces the prop unit with a marine jet motor using a pump to issue a jet of water to propel a boat. The Parker U.S. Pat., No. 5,356,319, is for a boat with a removably inboard jet propulsion unit in which the integral jet power unit is encased in a waterproof housing and positioned in a well located in the hull and is mounted to be removed from the hull.\nThe present invention is directed towards an outboard jet boat in which the main fuel tank and controls are mounted within the hull of a boat while the outboard jet drive unit is mounted in a housing with an engine and is removably attached to the transom of the boat. The fuel tank and controls are connected between the hull and outboard drive through quick disconnect couplings. The housing is shaped to support an engine on a platform directly over the jet drive unit for actuating the jet drive unit through a clutch mechanism with the engine and jet drive positioned parallel to each other."}
{"text": "Radar level gauge (RLG) systems are in wide use for determining the filling level of a product contained in a tank. Radar level gauging is generally performed either by means of non-contact measurement, whereby electromagnetic signals are radiated towards the product contained in the tank, or by means of contact measurement, often referred to as guided wave radar (GWR), whereby electromagnetic signals are guided towards and into the product by a probe acting as a waveguide.\nThe transmitted electromagnetic signals are reflected at the surface of the product, and the reflected signals are received by a receiver or transceiver comprised in the radar level gauge system. Based on the transmitted and reflected signals, the distance to the surface of the product can be determined.\nMore particularly, the distance to the surface of the product is generally determined based on the time between transmission of an electromagnetic signal and reception of the reflection thereof in the interface between the atmosphere in the tank and the product contained therein. In order to determine the actual filling level of the product, the distance from a reference position to the surface is determined based on the above-mentioned time (the so-called time-of-flight) and the propagation velocity of the electromagnetic signals.\nMost radar level gauge systems on the market today are either so-called pulsed radar level gauge systems that determine the distance to the surface of the product contained in the tank based on the difference in time between transmission of a pulse and reception of its reflection at the surface of the product, or systems that determine the distance to the surface based on the phase difference between a transmitted frequency-modulated signal and its reflection at the surface. The latter type of systems is generally referred to as being of the FMCW (Frequency Modulated Continuous Wave) type.\nMany products, in particular liquids, are at least to some extent transparent to the electromagnetic signals in the frequency range typically used in radar level gauge systems, which may possibly cause electromagnetic signals reflected from the bottom of the tank to interfere with the electromagnetic signals reflected from the surface of the product when measuring low filling levels, i.e. the distance resolution may be insufficient to distinguish the electromagnetic signals reflected from the bottom of the tank from those reflected from the surface of the liquid.\nEP 1 242 794 discloses a radar level measuring device having an absorber for microwave signals positioned at the base of the tank. Through this absorber, the reflection at the bottom of the tank can be considerably reduced, which allows measurement of low filling levels. In the radar level measurement device according to EP 1 242 794, the absorber is provided in the form of a connected plate of spikes where the tips of the spikes point in the direction of the incident microwaves. The spikes are made of a material which has a good absorption capacity for microwaves. Although the absorber according to EP 1 242 794 allows measurement of lower filling levels than was previously possible, the extension from the bottom of the tank to the top of the absorber limits the lowest filling level that can be measured to in the order of at least one wavelength of the microwave signals. As the frequency range of radar level gauging systems is typically in the order of tenths of gigahertz, the wavelength is in the order of a few centimeters."}
{"text": "Venous, arterial, and body fluid catheters are commonly used by physicians. For example, such catheters may be used to gain access to the vascular system for dialysis, for introducing pharmaceutical agents, for nutrition or fluids, for hemodynamic monitoring, and for blood draws. Alternatively, catheters can be used for drainage of fluid collections and to treat infection. Following introduction into the patient, the catheter is secured to the patient. In conventional practice, the catheter is commonly secured to the patient using an adhesive tape patch or by suturing an attached hub to the patient's skin."}
{"text": "It is known to introduce elongated elements into ducts with the help of compressed air. This method, called “jetting”, is described in the document U.S. Pat. No. 5,197,715 as well as the use of a “sonic head”, where “jetting” is a synergy of pushing and blowing and where at the “sonic head” a sonic flow is created, generating a pulling force at the foremost end of the cable. Blowing creates a distributed forward propelling force onto the cable, compensating locally the friction forces, limiting build-up of axial forces in the cable, in turn limiting extra exponential friction build-up by the well known “capstan effect” (See e.g. http://en.wikipedia.org/wiki/Capstan_equation). As air is a compressible fluid, it expands when the pressure decreases towards the end of the duct, resulting the air propelling forces to being relatively small in the first section of the duct (injection side) and relatively large at the last section of the duct (exhaust side). During jetting, typically there is a deficiency of air propelling forces to compensate the friction in the first section, where pushing assists the cable installation, and an excess of air propelling forces in said last section. The so-called “critical point” is the location in the duct where the air propelling forces have just become equal to the friction forces. In a situation where the cable weight is dominant for the friction (as is the case in jetting), the friction forces are mainly constant, and it is easily understood that the “critical point” is the most difficult point to reach: upstream of this point there is less length for the pushing force to build up and downstream of this point a pulling force will build up first. It shall be noted that when the cable diameter becomes close to that of the internal diameter of the duct, this reasoning does not hold anymore, because the pressure drop over the part of duct that is filled with cable will be higher because of the increased hydraulic resistance, as the hydraulic diameter is significantly reduced by the presence of the cable. For this jetting method the “sonic head” can be advantageously used. Calculation in the document U.S. Pat. No. 5,197,715 showed that the improvement is enhanced because the “sonic head” is as it were “pulling” the cable over the “critical point”.\nWhen the duct is bended with a continuous radius over a great length, the aforementioned synergy between pushing and blowing disappears: the action of axial forces in the cable is much shorter here because the afore mentioned capstan effect is much larger. In this case, the air has to do the job all by itself, and installation will only go for duct lengths short enough such that the blowing generates sufficient air propelling forces from the very beginning, which will cause the installation length to decrease by about 50%. However, for today's larger relative cable diameters the increased hydraulic resistance of the part of the duct filled with cable, helps to “pull” the cable through the most difficult first section as in this first section, the hydraulic diameter is (significantly) reduced by the presence of the cable.\nAnother method to introduce an elongated element is known, with the help of water, where the liquid at least partly compensates the weight of the elongated element. This method is called floating, and the operating conditions are quite different from the jetting method, especially regarding the flow and speed of the liquid injected into the duct. Indeed, the speed of the liquid is quite low compared to compressed air, and there is no expansion along the pipe. In many cases, much longer cable installation lengths can be reached in straight ducts by floating than by jetting.\nHowever, it is still difficult to operate (i.e. to introduce or remove) elongated elements in a duct when the latter is bended with a continuous radius over a great length. A typical case is when the duct is coiled over a drum for a length, up to more than one thousand meters for some cases. Another typical case is when the duct is helically stranded with another duct or cable over its full length, for example for longer than one thousand meter. Jetting or floating always end up in a length being limited by the friction between the duct walls and the elongated element, created by the continuous bending. Calculations demonstrate that the jetting method often achieves a higher installation length than the floating method in such bended ducts, but some improvements are still needed to achieve higher and higher target installation lengths."}
{"text": "1. Field of the Invention\nThe present invention relates to a deskew circuit and a deskew method that corrects skews between clock and data, as well as among data on a data bus, in data transfers to transfer large volume data between packages at high-speed and to transfer clock and data in parallel.\n2. Related Background Art\nA conventional technology that transmits data at high-speed while taking data skew into consideration has been indicated. In such a technology, a delay amount for every bit in data on the receiving side is adjusted using a phase adjustment pattern. However, phases are matched to an external clock supplied to both the transmitting and receiving sides, and the technology therefore is not a method to transfer clock and data in parallel.\nThe conventional technology involves a method of adjusting only the clock in devices such as disk array control devices in which there is a large skew in a backplane due to the large number of packages and in which a high-speed transfer must be implemented in a data transfer method that transfers clock and data (bus data) in parallel in order to reduce power consumption. Although this method is effective when the transfer speed is slow and data pulse width is wide, it is difficult to correct the clock to a position that allows all data to be read correctly when the transfer speed is fast and the data pulse width is narrow. This becomes even more difficult when taking into consideration the junction temperature fluctuations of an LSI that has a built-in circuit that makes such an adjustment, power source voltage fluctuations, flip-flop set-up time and hold time. Further, the skew is small and the clock can be corrected when the transmission distance is short (e.g., approximately 10 cm–20 cm), such as when sending and receiving within the same substrate (package). However, the skew is larger and correcting the clock becomes very difficult when the transmission distance is long (e.g., approximately dozens to 100 cm), such as in a backplane. Skew adjustment becomes difficult when the transmission distance is long as in a backplane, and this also sets a limit to the transfer speed and makes high-speed transfer impossible."}
{"text": "The present invention relates to a closure and/or filling device for closing and/or filling a reservoir, and to a dispenser device for dispensing a fluid, which dispenser device includes such a reservoir. The present invention is particularly well suited to being used in a spray device for spraying a liquid, in particular in the pharmaceuticals field.\nIn medication sprays, and in particular in devices in which the pump is pre-filled during packaging of the device, as applies in particular to certain single-dose or two-dose devices, certain problems need to be solved.\nA first problem relates to the stability of the fluid and to contamination thereof. While it is stored, a liquid tends to evaporate and/or can be contaminated by external agents such as bacteria, or can even be oxidized by oxygen. To avoid any evaporation, it is necessary to close the reservoirs in airtight manner and to use impermeable materials for the component parts of the reservoirs To avoid any contamination, it is important to keep the number of different materials in contact with the fluid to as small a number as possible, and said materials should be inert relative to the fluid. Suitable materials for satisfying those two requirements are, for example, glass for the reservoir, and special rubbers that do not interact with the liquid for the sealing device(s).\nAnother problem relates to filling. The pharmaceuticals industry generally requires devices that are easy to fill and that are, if possible, capable of being assembled completely before they are filled with medication. Clearly, the handling of the medication is the most critical step as regards problems that can arise due to the surroundings, and in particular the risk of the fluid being contaminated. It is therefore important to have not only a suitable filling machine, but, above all, to protect the entire filling and assembly process effectively, which can be very complex and therefore also very costly. In addition, the possibility of having a device that is completely pre-assembled before it is filled is going to become increasingly important, in particular for the single-dose or two-dose sprays that are to serve in future use with influenza vaccines. In such a case, the fluid, e.g. the vaccine will be available only a few months (2 to 3 months) before it is used,-and millions of devices will need to be filled within a very short time. In that case, it would be very advantageous for it to be necessary merely to fill the reservoir of a device that has already been completely assembled, it being possible for such assembly to be performed at any time without it being necessary to work to such a short time constraint.\nAnother problem that arises relates to precompression. When sprays are used, the quality of the spray and the size of the particles depends on the pressure delivered to the liquid as it is being expelled. In order to make such devices less dependent on the actuating force delivered by the user, it is preferred to provide a system capable of delivering sufficient precompression to the liquid before it is expelled.\nCurrently-available technology does not make it possible to solve all of the problems, or to satisfy all of the requirements mentioned above, and only partial and complex solutions are available in the state of the art.\nThus, devices have been proposed in which a membrane of any shape, material, or type is fixed over the opening of the reservoir, said membrane being designed to be broken or torn by a suitable device when the spray is actuated. That solution makes it possible to solve the problem of the stability of the fluid contained in the receptacle, but it suffers from the drawbacks of being relatively costly because the operation of assembling the membrane onto the reservoir is quite complex. In addition, filling is necessarily performed before the reservoir is assembled in the device, and fixing the membrane requires substances, such as plastics materials, to bond or glue the membrane, and such substances can cause problems related to contaminating the fluid. In addition, no precompression is provided by that type of closure system.\nIn another solution disclosed by Document FR-2 684 304, when the device is actuated, a rubber stopper which also serves as a pump piston is pierced by a hollow needle which is connected to the spray nozzle of the device. That solution makes it possible to solve the problem of the stability and of the contamination of the fluid contained in the reservoir because it uses only glass and rubber, but it does not make it possible to overcome the problem related to filling the reservoir after it has been assembled. The reservoir must be filled separately before final assembly of the reservoir. Furthermore, no precompression is provided by that system because the rubber offers very little resistance to penetration by the needle. In addition, in some cases, the piercing provided by the needle can remove small particles of rubber that can obstruct the spray or that can be ejected into the user\"\"s nasal cavity. Moreover, although the device is relatively simple to close, it is complex to open, it requires a plurality of components, and is difficult to handle, in particular because of the presence of a needle.\nDocument U.S. Pat. No. 3,841,329 discloses a syringe provided with a stopper suitable for forming the dispensing piston when the syringe is actuated. That piston-cum-stopper is provided with a slit that opens when the syringe is actuated. No precompression is provided since the slit acts merely as a valve.\nAn object of the present invention is to provide a closure and/or filling system for a reservoir, in particular serving to be used in a spray device, which system does not suffer from the above-mentioned drawbacks.\nThus, an object of the invention is to provide a reservoir that has a simple and effective closure system which is opened by the pressure that is generated inside the reservoir when it is actuated, without needing a complex opening system that removes or pierces or withdraws portions of components.\nAn object of the present invention is also to provide such a reservoir that is simple and inexpensive to manufacture and to assemble.\nAn object of the present invention is also to provide a closure and/or filling system that makes it possible to use as few materials as possible in contact with the fluid contained in the reservoir, and in particular glass for the reservoir, for the closure element, rubber characterized by being inert relative to the fluid.\nAn object of the present invention is also to provide a closure and filling system for a reservoir, which system can be used as a liquid inlet for easy and fast filling, it being possible, in particular, for filling to be performed prior to complete final assembly of the fluid dispenser device that incorporates the reservoir.\nAn object of the present invention is also to provide a closure and/or filling system that opens to enable the liquid to be expelled only if a predetermined and sufficiently high pressure is generated when actuating a reservoir, so that the liquid is expelled with a certain amount of precompression. This is particularly advantageous for a spray system.\nAn object of the present invention is also to provide a reservoir provided with a closure and filling device and that satisfies all of the above requirements simultaneously, i.e. it guarantees the stability of the fluid contained in the reservoir, it prevents the fluid from being contaminated by chemical interaction, it enables the reservoir to be filled quickly and simply, in particular after final assembly of the device, and it guarantees that the fluid is expelled with precompression that can be predetermined.\nThe present invention thus provides a reservoir for a fluid spray, said reservoir containing a fluid to be dispensed, and provided with a closure and/or filling device for closing and/or filling said reservoir, said device comprising at least one closure element acting as a leaktight closure stopper when at rest, said at least one closure element comprising:\nan annular portion whose outside diameter corresponds substantially to the inside diameter of the reservoir, said annular portion being inserted in leaktight manner in said reservoir; and\na central portion having an end wall which closes off said reservoir in leaktight manner when in the rest position;\nsaid reservoir being characterized in that said end wall is provided with at least one slit suitable for opening under the effect of at least a predetermined minimum pressure, so that the fluid is dispensed with predetermined precompression that is sufficient to provide a good spray of fluid, said predetermined pressure threshold being defined by the geometrical shape of said end wall and/or by the thickness of said end wall and/or by the characteristics of the material of said end wall and/or by the dimensions of said slit.\nAdvantageously, said central portion of said at least one closure element is provided with an axial bore that is closed at one end by said end wall which is made of an elastically-deformable material, and said slit opens by means of said end wall being deformed in response to at least a predetermined minimum pressure being exerted on it.\nAdvantageously, said closure element is suitable for receiving the end of an actuating member for actuating a dispenser device for dispensing fluid, said actuating member being provided with a fluid expulsion channel and said end of the actuating member being made to be concave so that, under the effect of a predetermined pressure generated inside the reservoir when said actuating member is actuated, said end wall of the closure element deforms towards the outside of the reservoir, into said concave end, thereby opening the slit, and dispensing the fluid contained in the reservoir through said expulsion channel and with precompression.\nAdvantageously, the actuating member is provided with a shoulder that co-operates with an abutment zone of said closure element so as to transmit the actuating force and so as to move said closure element in said reservoir for the purpose of dispensing the fluid in the reservoir, with precompression.\nIn a variant embodiment, at least said closure element is suitable for receiving a filling fitting for filling the reservoir with the fluid before or after it is assembled in a dispenser device for dispensing fluid.\nAdvantageously, said filling fitting is provided with a filling tube which, by exerting a certain amount of pressure on said end wall, deforms it towards the inside of the reservoir so as to pass through said slit as open.\nAdvantageously, said filling fitting is provided with a radial flange co-operating with the closure element and with the wall of the reservoir so that the pressure exerted by the tube to enable it to pass through the slit does not displace said closure element inside the reservoir.\nIn a preferred variant embodiment, the closure and/or filling device comprises a first closure element acting as a leaktight closure stopper when in the rest position and as a piston while the fluid is being dispensed, and a second closure element acting as a leaktight closure stopper and/or as a filling element.\nAdvantageously, a hollow rigid member is disposed in said bore in said second closure element and is suitable for receiving a filling fitting for filling the reservoir, which fitting is provided with a filling tube which, by exerting a certain amount of pressure on said end wall, deforms said end wall towards the inside of the reservoir so as to pass through said slit as open.\nAdvantageously, said hollow rigid member is provided with a radial flange co-operating with the second closure element and with the wall of the reservoir so that the pressure exerted by the tube to pass through the slit does not displace said second closure element inside the reservoir.\nIn another variant embodiment, the reservoir is further provided with a third closure element disposed in * fixed manner between the first and second closure elements, said third closure element separating the reservoir in leaktight manner into two sub-reservoirs, each sub-reservoir containing a different fluid, e.g. a powder and a liquid, and said second closure element further acting as a piston to mix the two fluids in one of the sub-reservoirs.\nAdvantageously, the third closure element is provided with a hollow rigid member which is inserted in its central bore, which makes it possible to open the slit in the end wall under the effect of at least a minimum pressure generated in one of the sub-reservoirs, and which prevents such opening under the effect of a pressure generated in the other sub-reservoir.\nPreferably, when the filling tube of the filling fitting passes through said slit, it opens said slit so that an opening is formed around said tube, thereby enabling air to escape as the reservoir fills with the fluid.\nThe present invention also provides a fluid spray, said fluid spray comprising an actuating member and a reservoir.\nPreferably, the reservoir is made of glass, and the closure element(s) is/are made of a substance that is inert relative to the fluid contained in the reservoir, and said fluid is sprayed finely and with precompression each time the device is actuated."}
{"text": "The present invention relates to an air-distributor box for a heating or air-conditioning system of motor vehicles, and more particularly, to an air-distribution box with a first air-outflow orifice for connection of a lateral nozzle arranged laterally in a motor vehicle, a second air-outflow orifice arranged below the first air-outflow orifice for ventilating a motor vehicle foot space, and an air flap controlling an air stream to the second air-outflow orifice.\nIn the air-distributor box described, for example, in DE 3,925,726 A1, the arrangement of the second air-outflow orifice for ventilating the foot space underneath the first air-outflow orifice for the lateral nozzle presents the problem that, with the air flap opened, air is supplied to the foot space to only an insufficient extent when the air supply to the lateral nozzles is not throttled on the outside at the lateral nozzles by flaps present there. This insufficient air supply to the foot space becomes noticeable particularly when the heating is in operation.\nAlthough such insufficient ventilation of the foot space could possibly be rectified by making the orifice cross-section of the first air-outflow orifice to the lateral nozzle substantially smaller than the orifice cross-section of the second air-outflow orifice for ventilating the foot space, such a reduction in cross section is nevertheless forbidden on account of the air requirement of the lateral nozzle in the event of maximum ventilation or cooling of the passenger compartment.\nAn object on which the present invention is based is to provide an air-distributor box in which there is no need for the outer throttling of the lateral nozzle for the purpose of a sufficient heating of the foot space when the air flap is open.\nThe foregoing object has been achieved according to the present invention by providing that the first air-outflow orifice is associated with a throttle diaphragm arranged to control an orifice cross-section of the first air-outflow orifice, and the throttle diaphragm and air flap are positively coupled such that the throttle diaphragm maximally throttles the orifice cross-section of the first air-outflow orifice in the open position of the air flap allowing a full air stream to the second air-outflow orifice and exposes the orifice cross section of the first air outflow orifice completely in a closing position of the air flap shutting off the air stream to the second air-outflow orifice.\nAn advantage of the air-distributor box according to the present invention is that, as a result of the throttle diaphragm at the air-outflow orifice to the lateral nozzle and its coupling to the air flap for ventilating the foot space, when the air flap is opened for ventilating the foot space, the orifice cross-section of the air outflow orifice to the lateral nozzle is automatically throttled correspondingly. Thus, a sufficient ventilation of the foot space is ensured, without the need for additional manipulations on the lateral nozzle. When the air flap is closed, the throttle diaphragm exposes the first air-outflow orifice completely, and when it is opened completely the orifice cross-section of the first air-outflow orifice is reduced to the maximum extent. Between these two end positions, i.e., \"foot-space air flap closed\" and \"foot-space air flap fully open\", a continuous closing of the air-outflow orifice by the throttle diaphragm is obtained for the supply of air to the lateral nozzle.\nIn one preferred embodiment of the present invention, the throttle diaphragm is configured so that maximum throttling does not create a total air shut-off relative to the lateral nozzle, but there is still a residual cross section for a minimum air flow in the air-outflow orifice. This residual cross-section is preferably selected so that the residual air stream corresponds approximately to that which would flow out via the lateral nozzle in the event of an outer throttling of the lateral nozzle to less than half.\nA simple constructive implementation of the throttle diaphragm is achieved by using as a throttle diaphragm a plate-shaped rotary slide which is rotatable parallel to the plane of the first air-outflow orifice about a pivot axis oriented at right angles thereto. The rotary slide preferably has the form of a circular segment, and its pivot axis is arranged near one edge of the first air-outflow orifice of rectangular configuration. If the first air-passage orifice is, for example, made square and the form of a segment of a circle of the rotary slide has a segment angle of approximately 90.degree. and a circle radius which is a little smaller than the side length of the square air-outflow orifice, then approximately the above-defined maximum throttling of the air passage cross-section in the first air-outflow orifice is achieved.\nA coupling rod invention is used for the positive coupling of the throttle diaphragm and air flap, which coupling rod, with the pivot axes of the throttle diaphragm and air flap oriented in parallel, is articulated on the air flap and on the throttle diaphragm of joint axes parallel to the pivot axes."}
{"text": "In recent years, a holey fiber becomes popular as an optical fiber suitable for a long transmission distance. The holey fiber is an optical fiber in which a refractive index of a clad is reduced by holes. It is known that the holey fiber obtains an optical characteristic that cannot be obtained by a conventional optical fiber in which a refractive index of a clad is reduced by an impurity.\nAs described in Patent Literatures 1 through 4, the holey fiber is manufactured by the steps of (1) preparing a columnar base material (hereinafter, referred to as a “preform”) made from silica glass, (2) forming, in the preform, through holes which are to be holes, and (3) drawing the preform in which the through holes have been formed.\nThe optical characteristic of the holey fiber is influenced by positions of the holes. Accordingly, it is important to form the through holes in predetermined appropriate positions in the preform in order to manufacture an optical fiber having a desired optical characteristic.\nIn the step (3), through holes extending in a direction vertical to end surfaces of the preform are formed by a drilling process. If a machine tool for forming the through holes has low machining accuracy, perforating positions of the through holes are gradually shifted as the perforating is proceeded even if the perforating is started from appropriate positions of one end surface. Accordingly, positions of the through holes are largely shifted at the other end surface, or a through hole is connected to another through hole in the middle of the preform. For this reason, after the through holes are formed, it is necessary to whether or not through holes are formed in respective appropriate positions.\nThe inspection as to whether or not through holes are formed in respective appropriate positions has been normally carried out by a method in which a preform is observed from an end surface by use of an optical microscope.\nHowever, the conventional inspecting method by use of an optical microscope can inspect only forming positions of holes in the vicinity of an end surface of a preform. That is, the forming positions of the holes in the middle of end surfaces (at an arbitrary cross-section between both end surfaces) cannot be inspected. There is another method in which forming positions of through holes in the middle of a preform are estimated from (i) forming positions of holes in one end surface and (ii) forming positions of the holes in the other end surface. However, such estimation cannot accurately specify the forming positions of the holes in the middle of the end surfaces. In view of the circumstances, it is necessary to inspect which portion of a preform is suitable for manufacturing an optical fiber by cutting the preform in round slices in a case where forming positions of through holes in an end surface are largely shifted from respective appropriate positions. This inspection increases a manufacturing cost.\nOn the contrary, Patent Literature 5 discloses a method for inspecting whether or not through holes are formed in respective appropriate positions on the basis of an intensity distribution of forward scattered light which is generated by parallel light entering from a side surface of a preform."}
{"text": "1. Field of the Invention\nThe present invention relates to an information recording medium (or data structure), an information recording/playback method, and an information recording/playback apparatus, which are suited to record/play back a digital stream signal used in digital TV broadcast or the like.\n2. Description of the Related Art\nIn recent years, TV broadcast has entered the era of digital broadcasts having Hi-Vision programs (programs of high-definition AV information) as principal broadcast contents. The current digital broadcast (and forthcoming terrestrial digital broadcast) adopts an MPEG transport stream (to be abbreviated as MPEG-TS as needed hereinafter). In the field of digital broadcast using moving pictures, MPEG-TS will be used as a standard format in the future. At the start of such digital TV broadcast, market needs for a streamer that can directly record digital TV broadcast contents are increasing.\nAs an example of a streamer that utilizes an optical disc such as a DVD-RAM or the like, “recording/playback apparatus” (Jpn. Pat. Appln. KOKAI Publication No. 2002-84479) is known.\nDifferent digital broadcast schemes are adopted in respective countries: for example, DVB (Digital Video Broadcasting) in Europe; ATSC (Advanced Television Systems Committee) in U.S.A.; and ARIB (Association of Radio Industries and Businesses) in Japan.\nIn DVB, the video format is MPEG2, the resolutions are 1152*14401, 1080*1920(i, p), 1035*1920, 720*1280, (576, 480)*(720, 544, 480, 352), and (288, 240)*352, the frame frequencies are 30 Hz and 25 Hz, the audio format includes MPEG-1 audio and MPEG-2 Audio, and the sampling frequencies are 32 kHz, 44.1 kHz, and 48 kHz.\nIn ATSC, the video format is MPEG2, the resolutions are 1080*1920(i, p), 720*1280p, 480*704(i, p), and 480*640(i, p), the frame frequencies are 23.976 Hz, 24 Hz, 29.97 Hz, 30 Hz, 59.94 Hz, and 60 Hz, the audio format includes MPEG1 Audio Layer 1 & 2 (DirecTV) and AC3 Layer 1 & 2 (Primstar), and the sampling frequencies are 48 kHz, 44.1 kHz, and 32 kHz.\nIn ARIB, the video format is MPEG2, the resolutions are 1080i, 720p, 480i, and 480p, the frame rates are 29.97 Hz and 59.94 Hz, the audio format includes AAC (MPEG-2 Advanced Audio Coding), and the sampling frequencies are 48 kHz, 44.1 kHz, 32 kHz, 24 kHz, 22.05 kHz, and 16 kHz.\nIn this manner, streams to be decoded by the apparatus side have many different variations for respective areas, and if all these variations are supported, the arrangement of a recording/playback apparatus (e.g., a DVD recorder) becomes considerably heavy (or completed). As a result, the apparatus cost increases."}
{"text": "In general, military specifications define a fuze or fusing system as a physical system designed to sense a target or respond to one or more prescribed conditions such as elapsed time, pressure or command, and initiate a train of fire or detonation in munitions. Safety and arming are primary roles performed by a fuze to preclude ignition of the munitions before the desired position or time. For example, if a munitions is an array of energetic charges launched from a watercraft into shallow (or deep) water to clear mines and obstructions, the fuze should ideally determine that launch and deployment occurred as intended before initiating detonation. Several fuzes designed to just prevent in-air detonation are known in the prior art.\nU.S. Pat. Nos. 3,765,331 (the '331 patent) and 3,765,332 (the '332 patent) disclose water-armed air-safetied detonators in which a plurality of small explosives are aligned in a spaced-apart fashion in a fuze housing. The first of these small explosives is a delay charge which ignites when the ordnance is launched or released. The delay charge eventually burns and flashes down an adjacent housing bore to ignite a transfer charge. Detonation of the transfer charge releases energy that is either used to move a piston (the '331 patent) or is in the form of a shock wave (the '332 patent). This released energy is delivered to a chamber that is flooded with either air or water depending on the environment in which the fuze is immersed. Adjacent the flooded chamber is a firing pin/percussion primer (the '331 patent) or just a percussion primer (the '332 patent). If the flooded chamber is filled with air, the released energy in the form of a moving piston (the '331 patent) or shock wave (the '332 patent) will not transfer through to the next stage of the fuze. If, however, the fuze is submerged in water, the flooded chamber is filled with water and the released energy entering the flooded chamber is transferred therethrough to the next stage of the fuze.\nAlthough being air-safed, these devices still have several disadvantages. Use of stored energy (i.e., explosive material) for arming and firing is considered bad design practice because the energy is available at all times during storage and transportation, and may therefore be released due to unforeseen causes or situations. The use of explosives as part of the fuze train can be inherently problematic. These problems range from the safety concerns related to the construction and storage of such devices to the fact that these fuzes are not reusable.\nTo overcome the problems inherent with the use of explosives, a mechanical underwater firing mechanism is disclosed in U.S. Pat. No. 2,660,952. Briefly, a spring-loaded plunger is mounted in a housing. The head of the plunger is formed with a recess. Fitted in the housing coaxial with the plunger is a plug having a central bored portion in which a firing pin is temporarily positioned intermediately therein by a shear pin. As a result, small chambers are defined in the central bored portion on either side of the firing pin. The central bored portion of the plug opposes the plunger's recess and is sized at its exterior to fit within the recess. When the spring-loaded plunger is cocked, the head of the plunger is spaced apart from the central bored portion of the plug to define a chamber within the housing. An opening in the side of the housing at the chamber allows the environment surrounding the housing (e.g., air or water) to fill the chamber.\nWhen the spring-loaded plunger is released, the plunger recess envelops the central bored portion of the plug to compress any fluid trapped in the small chamber of the central bored portion between the firing pin and the head of the plunger. If the trapped fluid is air, the compression thereof will not develop forces sufficient to cause the shear pin to fail. However, if the trapped fluid is water, the compression forces imparted by the plunger will be sufficient to cause failure of the shear pin thereby allowing movement of the firing pin to impact a primer.\nWhile eliminating the use of explosives in the firing mechanism, this device has other disadvantages. For example, the requirement that a small chamber be defined in the central bored portion opposing the head of the plunger raises the possibility that an air bubble will form therein when the device is submerged in water. The presence of such an air bubble could prevent the mechanism from functioning underwater. At the same time, the requirement that the small chamber be present in the central bored portion could also bring about an unwanted firing. This could occur if the mechanism were not cocked and inadvertently dropped in water. Water could seep into the mechanism and fill the small chamber. Then, an in-air release of the (cocked) plunger could bring about movement of the firing pin just as if the mechanism were submerged in water. Another disadvantage brought about by the requirement of the small chamber arises in sub-freezing environments. Specifically, water in the small chamber could quickly freeze due to its small volume. If this occurs, the mechanism will not function."}
{"text": "1. Technical Field\nThe present invention relates to a stacking mechanism and a recording apparatus.\n2. Related Art\nRecording apparatuses developed to date that record data on a recording medium include those of a type that execute ink jet printing with a recording head, and of another type that form a static latent image on a photosensitive material and transfer a toner image to the recording medium, to thereby record the image thereon. The latter can be found, for example, in JP-A-2005-348339.\nNow, the recording apparatus includes a stacking mechanism in which the recording medium, typically paper in sheet form, is to be stacked, and a user takes out the paper that has undergone the printing process from a paper output tray of the stacking mechanism, and refills the paper in a paper feed tray of the stacking mechanism.\nIn some of lately developed home-use printers, the paper output tray is superposed on the paper feed tray inside the outlet of the paper, in response to the requirement for reduction in size of the printer. To refill the paper in the printer thus designed, the user has to draw out the paper feed tray through the outlet.\nThe structure in which the paper output tray is superposed on the paper feed tray has, however, a drawback that the accessibility to a pull portion of the paper feed tray is degraded. This leads to inconvenience in refilling the paper, and a solution therefor is being sought."}
{"text": "The invention relates to a clamping chuck for holding clamped a brush body, in particular a toothbrush body in its head region, comprising at least two clamping jaws spaced apart from one another in the longitudinal direction of the brush body to be held.\nBrushes, in particular toothbrushes, are fed successively during production to a plurality of processing stations, for example a tufting device, a shearing device and a grinding device. Here, the brush bodies can be transported on a conveyor belt or a transport chain, with the individual brush bodies each being held in a clamping chuck.\nIt is desirable here for as many clamping chucks as possible to be arranged next to one another in the narrowest possible space so that the cycle width can be reduced and the switching times can be reduced and the processing of more brush bodies per unit time is possible. In addition, the wear must be reduced and it is possible for more brushes to be processed simultaneously on a grinding disk with a certain diameter if the clamping chucks in which the brush bodies are held are arranged closer to one another.\nA problem, however, is that, in order to open the clamping chucks for inserting and removing a brush body, mechanical devices are required which require a certain minimum spacing between adjacent clamping chucks."}
{"text": "1. Technical Field\nThe invention relates generally to lighting mechanisms as used on vehicles for lighting ahead of and behind the vehicle, and for indicating braking, turning and/or backing up of the vehicle. More particularly, the invention relates to lighting mechanisms which function as tail lights, brake lights, turn signals and/or backing lights as used on trailers, towed vehicles or other tow behind apparatus. Specifically, the invention is a portable, radio controlled light bar including various lighting features such as tail lights, brake lights, turn signals and/or backing lights and is affixable to a trailer or towed vehicle, and a radio wave transmitter box affixable to the tow vehicle and electrically connectable to the wiring system and specifically to the tail lights, brake lights, turn signals, and/or backing light electrical systems.\n2. Background Information\nAs is well known in the art of vehicles such as cars and trucks, vehicles are pre-wired from the factory with tail lights, brake lights, turn signals and backing lights, all of which are generally required by local ordinances and state laws. Similarly, all trailers and other tow-behind apparatus that are used on any roadways must also have at least tail lights and brake lights, and often also turn signals although arm signals as used by the driver out the driver's window are generally sufficient during day light hours in some jurisdictions. For these reasons, substantially all vehicles and trailers are pre-wired and initially sold with all of the necessary lights thereon.\nMany vehicles and trailers are not used on roadways at all or are principally used off road such as on the farm, construction sites or other off road area. However, occasionally these trailers must travel along roadways such as during moving of the farm trailer from farm to farm, or during moving of a construction trailer from job site to job site. In many instances and particularly with trailers used principally off road, these trailers are not sold with pre-wired lights. This is problematic when the trailers are to be towed on roadways as most if not all jurisdictions require at least certain minimum lights such as tail lights and brake lights.\nAs a result of this lack of lighting on many trailers, these trailers are either used illegally without lights, not driven on the roadways, or must have after market lighting systems installed. The towing of the trailer without lights is a dangerous event as significant property damage and personal injury may result from any mishap that may be caused due to lack of tail lights and/or brake lights including failure of another vehicle to see the trailer or see the trailer stopping. Such property damage and personal injury, particularly when resulting from a violation of the law, places the trailer owner and/or operator in a position of potentially high legal liability. The alternative of not driving the trailer on roadways is really not an option because generally it must be moved from site to site, and trailering of a trailer is a noneconomical way of moving the trailer.\nThe other alternative of \"after market\" wiring is generally the avenue chosen by such trailer owners. However, this is often expensive, time consuming, and sometimes rather difficult to install. Holes must often be drilled in the trailer and wires run underneath the trailer along the frame or body from the back of the trailer to the tongue. This option is also permanent in that each after market wiring system is permanently attached to that specific trailer, and for this reason an \"after market\" system must be purchased for every trailer.\nAnother towing area in which lighting is a problem is in the tow truck usage area where a broken down vehicle is being towed. The tow truck tail lights, brake lights, turn signals and backing lights are not attachable in any manner to the lights of the towed vehicle. For this reason, tow trucks often have oversized lights, overhead flashing lights, and even high or upwardly extended lights thereon to hopefully be visible around and particularly behind the vehicle in tow. Generally the only precautionary measure available to tow truck towed vehicles is flashing hazard lights. Although these lights are often bright and eye catching, the se lights do not adequately replace actual brake lights, tail lights, turn signals, and backing lights as the hazard lights merely flash.\nA yet even further area where lighting is a problem is on cars towed behind motor homes and other campers. Generally, the cars have all of the necessary lights but these are not readily connectable to the motor home. The car must then be wired special for alternative use of the cars lights by the car when in use or the motor home when in tow. A similar situation is also often present when one moves its residence and must tow one of its cars behind a moving van. The only alternative at this time is to rent a tow trailer for use with the moving vehicle.\nObviously a lighting system is needed on all trailers that are used on roadways; however, many of these trailers are only infrequently used on roadways and thus a permanently installed system that is also often expensive and time consuming is not a preferred option. A lighting system is also needed for tow truck and tow vehicle combinations. Thus, a need exists for a lighting system that is temporarily affixable both mechanically and electrically to a trailer or towed vehicle while being readily removable for use on other trailers and towed vehicles. Specifically, a portable light assembly is needed that is interchangeable between multiple trailers and towed vehicles and thus more economical for use with seldomly road-towed trailers such as are used in the farming and construction industries, with broken down vehicles during towed transit, and with cars in tow behind motor vehicles. Additionally, the need requires such light assembly be lightweight, easy to attach and detach, easy to electrically connect, cost effective and simple to manufacture."}
{"text": "Zeolites are well known as typical examples of inorganic ion exchangers. Zeolites are capable of supporting various metals through cation exchange of Na+ ions or Ca.sup.2 + ions in the interior of their crystals for metal cations in an aqueous solution. Further zeolites have excellent characteristics with respect to specific surface area, heat resistance, water resistance, mechanical strength, etc. and are therefore widely used as gas adsorbing and separating agents, agents for treating waste water containing heavy metals, ion fixing agents and carriers for metal catalysts.\nZeolites are crystalline substances consisting primarily of SiO.sub.2 and Al.sub.2 O.sub.3, and the crystals thereof have three-dimensional structure and have regular minute pores. The composition is expressed generally by (M.sub.2, M')O. Al.sub.2 O.sub.3. mSiO.sub.2. nH.sub.2 O wherein M and M' are monovalent and bivalent metal ions, respectively, m is the coefficient of silica, and n is the coefficient of crystal water. Zeolites include the faujasite group (sodalite group), chabazite group and mordenite group which are different in the structure of crystals. The cations present in the voids or channels of zeolites can be exchanged for other metal ions, so that zeolites are utilized for softening hard water and separating off metal ions. When the alkali metal of zeolites is exchanged for bivalent or trivalent metal ions or hydrogen ions, the zeolite forms a strong solid acid, which is useful as an excellent catalyst for cracking petroleum and various carbonium ion reactions. Zeolites carrying silver, copper, zinc or the like supported thereon are kneaded with polymers to prepare antibacterial wrapping materials for preventing deterioration for use in the field of foods. Especially, the silver-incorporating zeolite is also excellent in ethylene-adsorbing ability and is therefore valuable for use.\nHowever, the zeolites incorporating silver, copper or zinc fail to fulfill the requirement of being inexpensive which is characteristic of inorganic materials because of material costs and complex manufacturing process, and are very expensive materials.\nAccordingly, the object to be achieved by the present invention is to develop a novel inorganic material which is usable in place of the silver-, copper- or zinc-incorporating zeolites and which has ethylene-adsorbing ability and antibacterial activity.\nThis object is fulfilled by exchanging Ca.sup.2 + ions, of these ions and monovalent cations, such as Na+ ions, which are present between the layers of a tobermorite for silver ions into the tobermorite although this has never been practiced in the past.\nA description will be given with reference to Na+ ion which is typical of monovalent cations.\nWe have directed attention to tobermorites which are fibrous layered compounds heretofore widely used as lightweight heat-insulating materials and noncombustible building materials because of their resistance to a high temperature of 650.degree. C. and excellent heat-insulating properties, and conceived the entirely novel idea that the Ca.sup.2 + ions, or these ions and Na+ ions which are present between the layers of the mineral can presumably be exchanged for silver ions by cation exchange. We have carried out intensive research based on this idea and consequently succeeded for the first time in the cation exchange of Ca.sup.2 + ions, or these ions and Na+ ions, which are present between the layers of a tobermorite, for silver ions. We have further found that this novel substance has ability to adsorb ethylene and antibacterial activity."}
{"text": "“Data Warehouses” are based on relational databases, and may include data management tools to extract data from various sources and manage that data. The data is transferred to the data warehouse in data packages that may include both data and metadata. Information related to the various data packages may be maintained in the data warehouse. All of the data and associated information is not necessarily held in the same database table. Rather, through the use of key fields, the data and associated information may be held in various tables that are linked together in an associated, relational database. The data warehouse operates in conjunction with, or “rides” on, the database.\nProblems can arise when the links, that relate various data package information to each other, are broken. This can happen in a number of scenarios, some of which may not be avoidable, such as, for example, if a server shuts down at a particularly inappropriate moment. Often the reason for the link failure will not be known.\nA link failure does not necessarily cause an immediate problem with operation of the data warehouse. However, if an operation on, or with, a data package (“a data request”) is attempted after a link failure, the data request will likely fail. One such data request involves an attempt to delete a data package. Since the links are broken, the system, in some operations, does not “see” the data package. This will cause the system to overlook the data package and not delete it from, for example, an entity that holds the data for reporting (“a data provider”).\nProblems associated with such “invisible” data packages, which can not be “seen” by the data warehouse system because of link failures, are call “zombie” problems. One such zombie problem relates to requests involving such data. Such requests are called zombie requests. Zombie requests can cause major problems for the data provider because the request gets “locked” within the data warehouse or the data provider, is never fulfilled, and can not be deleted. However, most significantly, further data can not be loaded until the request is removed. Therefore this problem needs to be resolved with dispatch.\nOne way of addressing the problem of zombie requests is to have the database administrators remove all of the data and then reload the data. Another way of addressing the problem is to manually identify the data associated with the link failure. These techniques are labor intensive, time consuming, costly, and don't always succeed."}
{"text": "The present invention pertains to a vehicle headrest, and particularly to one which is adjustable.\nThere exists a variety of proposed headrest constructions which allow a headrest to be adjusted to fit a particular driver or passenger. One adjustable headrest provides a pivoting adjustment as disclosed in U.S. Pat. No. 4,600,240 while a pneumatically operated bellows-type control is provided for another adjustable headrest as disclosed in U.S. Pat. No. 4,778,218. These headrests as well as others which include screw jack-type adjustments with scissor arms for supporting a movable plate with respect to a stationary plate typically are mounted to rods which extend into the seat back and provide vertical adjustment of the headrests.\nAlthough these systems provide the desired adjustment motion, some are either somewhat complicated to manufacture or require auxiliary equipment to operate and/or are relatively slow and cumbersome to use."}
{"text": "Humidification of air is an operation concerned with an interphase transfer of mass and energy that occurs when air is brought into contact with water in which the air is essentially insoluble. Depending on whether the water is in a form of a liquid or a vapor, there are two air humidification processes: a) An adiabatic process, in which the treated air is brought into a direct contact with water and the required evaporation heat is extracted from the air that is being humidified, and b) An isothermic process, in which a water vapor at atmospheric pressure is added to the treated air to increase its moisture, in which the heat energy of the treated air is unaffected.\nThe isothermic humidification process is usually carried out in a central air conditioning air duct system or in an open space, by distributing and mixing of a stream of atmospheric steam with a stream of air. The amount of steam that can be added to a stream of air is limited and depends on the dry bulb temperature and the absolute moisture content of the treated air. The steam for humidifying the air may be produced either at the location of the steam distributor in a compact humidifier, or it can be delivered to the steam distributor from a central steam boiler.\nTechnical and commercial literature indicate, that the current art compact isothermic humidifiers are producing the atmospheric steam in a sealed water tank by boiling and evaporating the incoming feed water at atmospheric pressure in a single stage evaporation process. The required heat is provided either by electric power via two or more electrodes or electric resistance heating elements submerged in the boiling water, or by a pressure steam delivered from a central steam boiler in a submerged heat exchanger.\nAn electric steam generating humidifier is disclosed in U.S. Pat. No. 4,239,956 issued Dec. 16, 1980 and reissued Oct. 30, 1990 as RE 33,414. Referred to in this patent are U.S. Pat. No. 3,386,659 and U.S. Pat. No. 3,436,967 as disclosing a steam generator in combination with duct work of forced air heating system.\nDisclosed in U.S. Pat. No. 4,564,746, issued Jan. 14, 1986 to B. W. Morton et. Al, is a cabinet type steam generating room humidifier.\nThe feed water used in compact isothermic humidifiers may be a city water, softened water, deionized water (DI) or reverse osmosis treated water (RO). Regardless of the feed water quality, all compact isothermic humidifiers are provided with a method to control the flow of the feed water into the water tank, a method to control the volume and the water level in the water tank, and a method to control the operating pressure in the water tank.\nAs the feed water is converted into steam, impurities which enter with feed water are concentrated in the water tank. Of concern are mainly the inorganic compounds of hard scale forming substances such as calcium and magnesium. Each substance has its own solubility limit in water solution. When its concentration exceeds the solubility limit, the excess substance precipitates and builds up a hard scale on the submerged electric resistance heating elements, electrodes, heat exchanger, and the water tank walls. The build up reduces the overall heat transfer rate. To maintain the performance and the efficiency of the humidifier, the water tank, the submerged heating elements and the heat exchanger are regularly cleaned, and the water tanks provided with the electrodes are regularly replaced at a considerable maintenance and replacement material cost.\nTo extend the operating period of the water tank, all isothermic humidifiers operating with feed water containing dissolved solids (DS) are provided with a method to control the concentration of total solids TS in the boiling water to reduce the hard scale build up rate. Most of them control the TS in the boiling water by a regular periodic blowdown of a mixture of feed water and boiling water which results in excessive consumption of feed water and with excessive heat loss with the blowdown water.\nThe steam generation process used in current art compact isothermic humidifiers involves steps including transfer of heat from the heat source into the water to cause it to boil, evaporation of the water, concentration of TS in the boiling water, separation of the produced steam from the boiling water, discharge of the produced steam from the water tank, replacement of the evaporated water with feed water, blowdown of boiling water to maintain an acceptable TS concentration in boiling water, and replacement of the blowdown water with feed water. All these process steps occur in a water tank at atmospheric pressure.\nIn current art compact isothermic humidifiers operating with feed water containing DS the flow of the feed water into the water tank is periodic and the steam generation process is cyclic and includes two operating periods.\nThe first operating period involves steps including transfer of heat from the heat source into the boiling water, evaporation of the boiling water, precipitation of DS and concentration of TS in the boiling water, separation of the produced steam from the boiling water, and discharge of the produced steam from the water tank.\nThe following second operating period, in addition to the above steps, involves process steps including filling up of the water tank with make up water, draining of a small portion of the boiling water, and heating of the feed water up to its boiling temperature. During the second operating period the concentration of DS and TS in the boiling water drops in proportion to the flow of the feed water. After a number of steam cycles the water tank is drained and flushed to minimize build up of the hard scale.\nBecause of the periodic intake of cool feed water into the water tank and use of a portion of the process heat in heating up of the feed water to boiling temperature, the rate of production of steam during the second operating period is reduced. Depending on the flow rate of feed water, temperature of the water in the water tank may fall below the boiling point. Under such conditions the steam is not produced and the treated air stream is not humidified.\nAll current art isothermic humidifiers use process controls to keep the process equipment and the integrated humidification and steam generation processes operating and carried out as required to maintain a preset humidity of the air being humidified. There are four groups of controls that are integrated to maintain the operation of the humidifying system fully automatic. They include: a) the humidification process controls regulating the humidification of a stream of treated air, b) the electric power and safety controls regulating the power input in response to a humidity load, or steam controls for regulating the steam input to the heat exchanger in response to a humidity load, c) the water controls regulating the flow of feed water, water level, concentration of TS, blowdown of the boiling water, and draining of the water tank, and d) the monitoring and display of selected performance parameters.\nA major concern with the current art compact isothermic humidifiers are the very high operation energy cost and operation maintenance cost. The high operation energy cost is the result of use of the electric power as the source of the heat required in the production of the steam. The high operation maintenance cost is due to the required regular cleaning of the water tank and of the submerged electric resistance heating elements or heat exchanger, or the regular replacement of the water tank operating with electrodes due to the excessive build up of hard scale.\nTherefore, to overcome the shortcomings of the current art humidifiers, the object of the present invention is to provide a compact apparatus that produces steam from feed water containing dissolved solids using a gaseous or a liquid fuel and at minimum capital, operating and maintenance costs."}
{"text": "In circuits for driving gas discharge lamp loads, such as fluorescent lamps, it is known to reduce the power from which the lamps are driven (so as to produce dimming of the lamps) by using a resonant inductor and capacitor in series with the lamps and by varying the circuit's operating frequency. In such a known circuit, when the operating frequency of the circuit is changed, current through the lamps is reduced and the lamps are correspondingly dimmed.\nHowever, employing variation of the circuit's operating frequency in order to produce dimming renders the actual dimming level of the circuit susceptible to changes in the circuit's temperature which cause the circuit's operating frequency to change."}
{"text": "Capacitance sensors that respond to charge, current, or voltage can be used to detect position or proximity (or motion or presence or any similar information), and are commonly used as input devices for computers, personal digital assistants (PDAs), media players, video game players, consumer electronics, cellular phones, payphones, point-of-sale terminals, automatic teller machines, kiosks and the like. Capacitive sensing techniques are used in user input buttons, slide controls, scroll rings, scroll strips and other types of sensors. One type of capacitance sensor used in such applications is the button-type sensor, which can be used to provide information about the existence or presence of an input. Another type of capacitance sensor used in such applications is the touchpad-type sensor, which can be used to provide information about an input such as the position, motion, and/or similar information along one axis (1-D sensor), two axes (2-D sensor), or more axes. Both the button-type and touchpad-type sensors can also optionally be configured to provide additional information such as some indication of the force, duration, or amount of capacitive coupling associated with the input. One example of a 2-D touchpad-type sensor that is based on capacitive sensing technologies is described in U.S. Pat. No. 5,880,411, which issued to Gillespie et al. on Mar. 9, 1999. Such sensors can be readily found, for example, in input devices of electronic systems including handheld and notebook-type computers.\nA user generally operates a capacitive input device by placing or moving one or more fingers, styli, and/or objects, near a sensing region of one or more sensors located on or in the input device. This creates a capacitive effect upon a carrier signal applied to the sensing region that can be detected and correlated to positional information (such as the position(s) or proximity or motion or presences or similar information) of the stimulus/stimuli with respect to the sensing region. This positional information can in turn be used to select, move, scroll, or manipulate any combination of text, graphics, cursors and highlighters, and/or any other indicator on a display screen. This positional information can also be used to enable the user to interact with an interface, such as to control volume, to adjust brightness, or to achieve any other purpose.\nAlthough capacitance sensors have been widely adopted for several years, sensor designers continue to look for ways to improve the sensors' functionality and effectiveness. In particular, engineers continually strive to simplify the design and implementation of position sensors without increasing costs. Moreover, as such sensors become increasingly in demand in various types of electronic devices, a need for a highly-flexible yet low cost and easy to implement sensor design arises. In particular, a need exists for a sensor design scheme that is flexible enough for a variety of implementations and powerful enough to provide accurate capacitance sensing while remaining cost effective.\nAccordingly, it is desirable to provide systems and methods for quickly, effectively and efficiently detecting a measurable capacitance. Moreover, it is desirable to create a design scheme that can be readily implemented using readily available components, such as standard ICs, microcontrollers, and discrete components. Other desirable features and characteristics will become apparent from the subsequent detailed description and the appended claims, taken in conjunction with the accompanying drawings and the foregoing technical field and background."}
{"text": "To effect infusion into a patient, an indwelling catheter connected with an infusion line is inserted into the blood vessel of the patient. This indwelling catheter is constructed of a hollow outer catheter and a hollow inner needle having a sloped edge face at its leading end that protrudes beyond the leading end of the outer catheter. The needle point of the inner needle is provided with an opening. The catheter is inserted into the blood vessel of the patient by first introducing the leading end of the needle into the blood vessel of the patient.\nWhen the inner needle reaches the inside of the blood vessel, the blood enters the opening at the needle point and flows or flashes back into a transparent housing which is mounted at the root portion of the inner needle. This makes it possible to confirm that the inner needle has reached the blood vessel. At this time, the inner needle and the outer catheter are moved forward a slight distance to insert the leading end of the outer catheter into the blood vessel.\nIn known indwelling catheter assemblies, the inner circumference of the leading end of the outer catheter is closely fitted to the outer circumference of the inner needle so that the blood will not flow between the outer catheter and the inner needle. This makes it impossible to visually confirm by flashback of the blood whether or not the outer catheter has captured or been inserted into the blood vessel.\nThus, after confirming capture of the blood vessel by the inner needle, the inner needle and the outer catheter are delicately advanced by the operator based on the senses and experience of the operator. Because of this, the proper insertion of the outer catheter into the blood vessel is not always achieved, but may be too shallow or deep. This is partly because the spacing between the edge face of the inner needle and the leading end of the outer catheter is slightly different depending upon the kind and size of the indwelling catheter assembly, and the manufacturer of the indwelling catheter assembly. This is also due in part to the fact that the size of blood vessels in a patient(s) varies.\nOne proposal that attempts to address this problem is described in U.S. Pat. No. 5,279,572. This patent discloses an indwelling intravenous needle that includes a metallic internal needle, an external needle positioned over the internal needle so that the leading end of the external needle is set back from the leading end of the internal needle, and a plastic housing fitted to the rear end of the exterior needle. The patent describes providing a longitudinal groove on the surface of the inner needle. This groove extends from the leading edge face of the inner needle and terminates at a point within the housing. During use of the indwelling intravenous needle, the leading end of the internal needle is inserted into the vein cavity and the visual observation of blood flashback within the internal needle confirms that the internal needle has captured the vein cavity. Further advancement of the indwelling intravenous needle eventually causes the leading end of the external needle to capture the vein cavity and this capture can be observed as a result of blood flashback within the groove in the surface of the internal needle. The groove allows blood to flow between the inner surface of the exterior needle and the exterior surface of the internal needle.\nThe indwelling intravenous needle disclosed in the aforementioned patent is, however, susceptible of certain disadvantages and drawbacks. In one respect, the patent describes forming the groove on the surface of the internal needle by machining. The nature of this machining process causes a significant reduction in the wall thickness of the internal needle, thus resulting in a severe weakening of the needle wall. This weakened wall can cause problems during use of the needle, particularly given the rather thin nature of the needle. Further, the machining process results in a groove having relatively sharp or rough corners where the sides of the groove meet the outer surface of the needle. These roughened corners can cause damage to the external needle. Also, because the front portion of the internal needle is exposed beyond the leading end of the external needle, these sharp corners of the groove present the potential for scraping against and possibly causing damage to the vein during insertion of the indwelling intravenous needle into the vein. It is possible to attempt to round these sharp edges through use of a secondary working operation such as a polishing or deburring operation, but this undesirably complicates the manufacturing process and increases the manufacturing time.\nIn light of the foregoing, a need exists for an indwelling catheter assembly which permits visual confirmation that the indwelling catheter has captured a blood vessel, but which does not present the possibility of damaging the surrounding portion of the assembly or needlessly injuring the vein."}
{"text": "The term “ionic liquid” (IL) is usually used to refer to a salt which is liquid at temperatures below 100° C., in particular at room temperature. Such liquid salts typically comprise organic cations and organic or inorganic anions, and are described inter alia in P. Wasserscheid et al., Angew. Chem., 2000, 112, 3926-3945.\nIonic liquids have a series of interesting properties: Usually, they are thermally stable, relatively non-flammable and have a low vapor pressure. They show good solvability for numerous organic and inorganic substances. In addition, ionic liquids have interesting electrochemical properties, for example electrical conductivity which is often accompanied by a high electrochemical stability.\nThese attributes give rise to many applications of ionic liquids: They can be used foremost as solvent in synthesis, as electrolyte, as lubricant and as hydraulic fluid. Moreover they serve as phase-transfer catalyst, as extraction medium, as heat-transfer medium, as surface-active substance, as plasticizer, as conductive salt, organic salt or additive in electrochemical cells, as electrolyte, as component in electrolyte formulations, wherein such electrolyte formulation comprising an ionic liquid is preferably used in electrochemical and/or optoelectronic device such as a photovoltaic cell, a light emitting device, an electrochromic or photo-electrochromic device, an electrochemical sensor and/or biosensor, particularly preferred in a dye sensitized solar cell.\nTherefore, there is a fundamental need for ionic liquids having a variety of properties which open up additional opportunities for their use.\nAn interesting family of ionic liquids contains tetravalent boron anions. Tetrafluoroborate containing ionic liquids were among the first of this new generation of compounds and 1-ethyl-3-methylimidazolium tetrafluoroborate ([EMIm][BF4]) was prepared via metathesis of [EMIm]I with Ag[BF4] in methanol as disclosed by J. S. Wilkes et al., J. Chem. Soc. Chem. Commun. 1990, 965.\nE. Bernhardt, Z. Anorg. Allg. Chem. 2003, 629, 677-685, discloses the reaction of M[BF4] (M=Li, K) with (CH3)3SiCN (TMSCN). The preparation of Li[BF(CN)3] is disclosed to take 7 days, that of K[BF(CN)3] takes one month. The yield of K[BF(CN)3] was 60%, the product contained 5% K[BF2(CN)2]. The molar ratio of [BF4]−:TMSCN was 1:7.8.\nUS 2011/150736 A1 discloses as a “Third Production Method” a reaction of three compounds: TMSCN, an amine or ammonium salt, and a boron compound.\nEP 2 327 707 A claims in claim 7 a method for producing an ionic compound represented by the general formula (I), comprising a step of reacting starting materials containing a cyanide and a boron compound. General formula (I) is a salt of a cation Ktm+ with [B(CN)4]−.\nThe examples disclose various methods for preparing tetrabutylammonium tetracyanoborate, for example:    1) Example 1-1 of EP 2 327 707 A discloses a reaction of tetrabutylammonium bromide, zinc (II) cyanide and boron tribromide in toluene at 130° C. for 2 days, with a yield of 35%. The molar ratio of boron compound:TMSCN was 1:5.5.    2) Example 2-1 of EP 2 327 707 A discloses a reaction of tetrabutylammonium bromide, tetrabutylammonium cyanide and boron tribromide in toluene at 130° C. for 2 days, with a yield of 77%. The molar ratio of boron compound:tetrabutylammonium cyanide was 1:7.1.    3) Example 3-3 of EP 2 327 707 A discloses a reaction of tetrabutylammonium bromide, trimethylsilyl cyanide and boron trichloride in p-xylene at 150° C. for 30 hours, with a yield of 98%. The molar ratio of boron compound:TMSCN was 1:5.5.    4) Example 3-11 of EP 2 327 707 A discloses a reaction of boron trifluoride diethyl ether, tetrabutylammonium bromide and trimethylsilylcyanide at 170° C. for 30 hours, with a yield of 75%.But not all embodiments which fall under claim 7 actually work well: Example 3 of the instant invention shows one embodiment also starting with boron trifluoride diethyl ether, which falls under claim 7, but produces the desired [B(CN)4] salt only as a by-product in negligible amounts, the main product is a [BF(CN)3] salt.\nThere was a need for a simplified method with high yield and satisfactory purity for the preparation of fluoro cyanide compounds of the 13th group of the periodic table with the anion having the general formula [(Z1F4−m(CN)m)−] with Z1 is B, Al, Ga, In or Tl and m being 1, 2, 3 or 4. The boron source should be a readily available compound with low costs. The cyanide source should not be a metal cyanide to avoid its negative impact on the environment. The number of reactants should be small and the method should allow the conversion without the presence of a solvent. The content of Cl and Br in the final product should be low. Also the content of Si and cyanide in the final product should be low. The method should require as few steps as possible. The method should allow also the preparation of compounds with m being 1, 2, 3 or 4 and not only of either a compound with m being 3 or a compound with m being 4. The method should avoid the use of Cl2, AgCN or AgBF4. The method should provide stable compounds of said formula which can be used as ionic liquids or as precursors of ionic liquids and can be used e.g. in electrolyte formulations and in electrochemical or optoelectronic devices. These compounds should be able to be disposed of in an environmentally friendly manner after use.\nThe method should allow the preparation of the desired compounds in high yields and under mild conditions with respect to methods disclosed in the prior art.\nThis object is achieved by a method using trimethylsilylcyanide as CN source and by doing the reaction in the presence of a Lewis acid. No Cl2, AgCN or AgBF4 is required. The content of Cl, Br, Si and cyanide in the final product is low. Another advantage is that the reaction does not require an extra solvent. The method has a reduced number of steps compared to the methods known from the prior art. The method allows for the preparation not only of compounds with m being only 3 or only 4, but for compounds with n being 1, 2, 3 or 4. These compounds can be prepared specifically and individually, and not only as mixtures. The reaction can be done under milder conditions than those used in the methods of the prior art, the reaction can be done at lower temperature or in shorter time.\nIn this text, the following meanings are used, if not otherwise stated:                alkyl linear or branched alkyl;        C1-q alkyl refers to any alkyl residue which contains from 1 to q carbon atoms; for example C1-6 alkyl encompasses inter alia methyl, ethyl, propyl, isopropyl, n-butyl, isobutyl, sec-butyl, tert-butyl, n-pentyl, isopentyl (3-methylbutyl), neopentyl (2,2-dimethylpropyl), n-hexyl and isohexyl (4-methylpentyl);        C2-q alkenyl refers to an alkenyl residue which contains from 2 to q carbon atoms and contains at least one double bond, the carbon chain can be linear or branched; for example C2-4 alkenyl encompasses inter alia ethenyl, 1-methylethenyl, prop-1-enyl, prop-2-enyl, 2-methylprop-2-enyl and buta-1,3-dienyl;        C2-q alkynyl refers to an alkynyl residue which contains from 2 to q carbon atoms and contains at least one triple bond, the carbon chain can be linear or branched; for example C2-4 alkynyl encompasses inter alia ethynyl, prop-1-ynyl and prop-2-ynyl;        C6-10 aryl refers to an aryl residue which has from 6 to 10 carbon atoms and is unsubstituted or substituted by 1, 2, 3 or 4 identical or different substituents independently from each other selected from the group consisting of C1-4 alkyl and C1-4 alkoxy; for example C6-10 aryl encompasses inter alia phenyl, methylphenyl, methoxyphenyl, dimethylphenyl, ethylmethylphenyl, diethylphenyl and naphthyl;        cyclic alkyl or cycloalkyl include cyclo and polycyclo, such as bicyclo or tricyclo, aliphatic residues;        C3-q cycloalkyl refers to a cycloalkyl group having from 3 to q carbon atoms; for example C3-10 cycloalkyl encompasses inter alia cyclopropyl, cyclobutyl, cyclopentyl, cyclohexyl, cycloheptyl, cyclooctyl, cyclononyl and cyclodecyl;        C1-q alkoxy refers to an linear or branched alkoxy group having from 1 to q carbon atoms; for example C1-20 alkoxy encompasses inter alia methoxy, ethoxy, propoxy, isopropoxy, n-butoxy, isobutoxy, sec-butoxy, tert-butoxy, pentyloxy, 1,4-dimethylpentyloxy, hexyloxy, heptyloxy, octyloxy, 1,5-dimethylhexyloxy, nonyloxy, decyloxy, 4-ethyl-1,5-dimethylhexyloxy, undecyloxy, dodecyloxy, tridecyloxy, tetradecyloxy and eicosyloxy;        alkylene means a linear or branched alkylene group; e.g. propylene, and e.g.        \npropylene can be connected via its C1 and C2 carbon atoms (a branched alkylene group), or via its C1 and C3 carbon atoms (linear alkylene group);                BMMIm n-Butyl-2-methyl-3-methylimidazolium        \n                BMPy n-Butylmethylpyridinium        \n                BMPyrr n-Butylmethylpyrrolidinium        \n                BMPip n-Butylmethylpiperidinium        \n                DCM dichloromethane;        EMIm 1-ethyl-3-methylimidazolium        \n                eq. molar equivalent;        halide F−, Cl−, Br− or I−, preferably F−, Cl− or Br−, more preferably Cl−;        halogen F, Cl, Br or I; preferably F, Cl or Br;        HEIm 1-ethylimidazolium        \n                IL ionic liquid;        “linear” and “n-” are used synonymously with respect to the respective isomers of alkanes;        RT room temperature, it is used synonymously with the expression ambient temperature;        Tdec decomposition temperature;        THF tetrahydrofuran;        TMSCN (CH3)3SiCN, i.e. trimethylsilylcyanide;        Trityl means the trityl cation, i.e. [Ph3C+]        “wt %”, “% by weight” and “weight-%” are used synonymously and mean percent by weight.        The expressions dye sensitized solar cell and photosensitized solar cell are used synonymously.        "}
{"text": "Barcode symbols provide a fast and accurate means of representing information about an object. Decoding or reading a barcode is accomplished by translating the pattern of barcode features, such as bars and spaces in linear (1D) barcodes or blocks or other features in a two-dimensional (2D) barcode, into the corresponding numbers or characters. Barcodes are widely used for encoding information and tracking purposes in retail, shipping and industrial settings.\nBarcodes are typically read by a handheld barcode reading device that includes a camera or imager (these terms being used interchangeably herein) for capturing an image of the barcode. Since the reading device is handheld, the quality of the captured image can suffer due to hand movement of the reading device relative to the barcode during the image capture period. For example, movement in the direction perpendicular to the bars and spaces with a linear barcode can result in a blurring that causes misdetection of the edges between the bars and spaces in decoding the barcode. For 2D barcodes, motion in any direction can induce image blur that results in misdetection of the barcode feature. The image blurring problem can be exasperated when using low cost image sensors that require longer exposure periods to capture images with suitable signal to noise ratios or when attempting to capture an image of a barcode under less than ideal lighting condition which also requires a relatively longer exposure period.\nOptical stabilization techniques are often used to reduce blurring associated with the motion of a camera during exposure to compensate for movement of the imaging device. Optical stabilization varies the optical path to the image sensor of the camera in response to movement. This technology can be implemented in the lens itself, or by moving the sensor as the final element in the optical path. The key element of all optical stabilization systems is that they stabilize the image projected on the sensor before the sensor converts the image into digital information. Optical stabilization techniques typically require additional sensors, such as gyroscopes and accelerometers, and devices to move lenses or sensors in the optical path. These additional sensors and devices create a number of potential disadvantages in a handheld device, such as increasing the expense of the handheld device, increasing the number of moving parts in the handheld device, possibly reducing the overall reliability, and increasing the power consumption within the handheld device.\nAnother known technique used to minimize motion blur relies on using an accelerometer to detect pauses in hand motion. This technique only allows the camera to attempt to capture an image during a detected pause in motion. This technique requires knowledge of the expected motion pattern of hand jitter in order to correctly detect the pauses. This technique is used by some cell phone cameras, one example being the Glogger VS2 software available for Nokia phones running on Symbian OS. Published U.S. Patent Application 20050236488 to Kricorissian, and assigned to the assignee of the present invention, discusses another implementation of a delay-based system. These approaches have proven not to be practicable for barcode readers because they often create an unacceptable delay when capturing the image of the barcode. Furthermore, these approaches may still result in blurring of the captured image in lower light environments or when using a low-cost image sensor where a longer exposure time is required since these approaches do not prevent blurring or compensate for blur causing movement, but instead merely attempt to avoid capturing the image of the barcode when movement is occurring.\nAn alternative approach to reduce image blurring is to reduce exposure times. However, imagers used in barcode reading devices are typically lower cost image sensors that require relatively longer exposure times. In the scanning environments where barcode reading devices are typically used, the light levels are often not sufficient to use reduced exposure times to eliminate image blurring. Providing supplementary lighting, such as a flash or other light source, can reduce the exposure time, but would result in increased power consumption by the handheld device with a commensurate reduction in the useful operating time of the battery powered handheld device."}
{"text": "This invention relates generally to dynamoelectric machines such as large turbo-generators and more particularly to cooling the end turns of the rotor winding in such machines.\nA rotor for a dynamoelectric machine typically comprises a cylindrical forging of magnetic metal having a plurality of axially-extending slots formed therein at circumferentially spaced positions. Conductor bars are disposed in the slots for carrying current. The ends of the conductor bars are suitably connected with conductive end turns to form the required current pattern. Because the conductor bars and end turns give rise to resistive heating, means for cooling the rotor winding are required. Effective cooling of the rotor winding contributes directly to the output capability of a dynamoelectric machine.\nRotor endwindings are difficult regions to cool effectively, as evidenced by the various schemes that have been devised for this purpose. Conventional techniques include providing cooling passages directly in the end turn conductors (\"direct cooling\") or creating regions of relatively higher and lower pressures to force cooling gas to pass over conductor surfaces (\"forced cooling\"). These conventional techniques add considerable complexity and cost to rotor construction. For example, directly cooled conductors must be machined or fabricated to form the cooling passages. In addition, an exit manifold must be provided to discharge the gas into the rotor. Forced cooling schemes require the rotor end region to be divided into separate pressure zones, with the addition of numerous baffles, flow channels and pumping elements.\nAn alternative to these expensive active cooling schemes is passive cooling, sometimes referred to as \"free convection cooling.\" Passive cooling relies on the centrifugal and rotational forces of the rotor to circulate gas in the blind, dead-end cavities formed between the concentric end turns. While passive cooling provides the advantages of minimum complexity and cost, heat removal capability is comparatively weak compared to active cooling because of low circulation velocities inside the cavities. One known approach to increasing gas velocity in the cavities is to machine holes in the retaining ring which encloses the end turns. However, because the retaining ring is under very high stress, openings therein can lead to mechanical failure.\nAccordingly, there is a need for a passive cooling system for cooling rotor endwindings with improved heat removal capability without diminishing the structural integrity of the retaining rings."}
{"text": "A flash memory is typically configured by an array of multiplicity of memory cells aligned in the word line direction and the bit line direction. A memory cell is configured by a stacked gate structure in which a floating gate electrode, interelectrode insulating film, and control gate electrode are stacked in the listed sequence. As flash memory increases its storage capacity through densification, features within the memory cell are packed in tighter dimensions. Dimensions typically affected by the densification are widths of floating gate electrodes and element isolation trenches. For instance, narrowing of the element isolation trenches makes formation of control gate electrode difficult since polysilicon film, typically employed as a control gate electrode material, needs to be filled in the narrowed gaps between the neighboring floating gate electrodes which are further narrowed by the presence of the interelectrode insulating film. One solution to this problem may be thinning the interelectrode insulating film. However, as the memory cells become smaller, relatively higher electric field is applied to the floating gate electrode through the interelectrode insulating film during programming. This is because downscaling of a memory cell often results in a floating gate electrode with a sharp tip. Because electric field tends to concentrate at the tips, edges, and corners, the downscaled floating gate having relatively larger percentage of such high electric field regions are subjected to larger amount of high electric field leakage current which prevents the memory cell from being programmed to the desired threshold.\nThe interelectrode insulating film is often configured by an ONO structure in which a silicon nitride film is interposed between the top silicon oxide film and the bottom silicon oxide film. The problem encountered in such ONO structure is electron traps within the silicon nitride film, which typically occurs in the aforementioned high-field regions located at the top portion of the floating gate electrode where electric field concentration occurs. The trapped electrons, when detrapped, degrades the charge retainability of the interelectrode insulating film.\nIt is further known that thinner interelectrode insulating film results in shorter endurance, in other words, reduced voltage tolerance. The interelectrode insulating film is not only used in the gate structures in the memory cell region but also in the elements for driving device operation provided in the peripheral circuit region such as capacitors for generating voltage. In applications where the device is expected to run extensively, enhanced endurance, i.e. voltage tolerance of the interelectrode insulating film is an important factor in achieving a long running device. However, a thinner interelectrode insulating film, which is advantageous in terms of gap fill capability, degrades the endurance/voltage tolerance of the interelectrode insulating film, and thus, will subject the elements in the peripheral circuit region and consequently the entire device to greater risk of malfunctioning. It is thus, required to further improve the insulativity of the interelectrode insulating film."}
{"text": "When visiting a theme park, such as Disney and Universal, for example, it is common for guests to be approached by a theme park photographer to have their photo taken as a souvenir that can later be retrieved at a designated area. For each park guest and their group having their photo taken, a printed ticket with a number corresponding to their particular photo is given to them by the theme park photographer.\nThe number on the printed ticket typically corresponds to a bar code that is associated with the photo or photos taken by the theme park photographer of the park guest and their group. A different bar code is used for each park guest and their group having their photo taken.\nTo associate the bar codes with the photos being taken, a personal digital assistant (PDA) is used. Each theme park photographer has a PDA coupled to their camera via an interface cord. The theme park photographer uses the PDA to first read or scan a bar code. The bar code is usually taken from a tablet of bar codes carried by the theme park photographer. After the photos of a park guest and their group have been taken, they are then transferred to the PDA. The PDA includes software that associates the bar code with the corresponding photos taken of the park guest. This process is repeated for each group of park guests having their photo taken.\nThe theme park photographer later transfers the stored photos and corresponding bar codes in the PDA to a photography database. Prior to leaving the theme park, the park guest can present their printed ticket at a photography gift shop. The photography database is searched using the number on the printed ticket which corresponds to the bar code used to identify their photos.\nA drawback of using theme park photographers to capture the experiences of theme park guests is cost. The use of theme park photographers is labor intensive. In addition, theme park photographers typically roam around the theme park, which means that park guests may miss out on the opportunity to have their photo taken a certain spots that are very scenic and reflective of their experiences at the theme park."}
{"text": "Corn cobs have, for many years, added a certain measure of value to the harvest. Prior to the early 60's, the common corn harvesting practice was to pick the ears in the field, transport the crop to corn cribs, and later shell the corn off the cob at stationary shellers. The need to dispose of the cobs, after shelling, spurred most of the early research for useful corn cob end products. However, using stationary shellers limited the volume of cobs which had to be dealt with.\nThe advent of in-field shelling, by self-propelled combine harvesters, created extreme efficiencies. The combine harvesters were designed to recover only shelled corn, or shelled corn incorporating only small amounts of cob, while crushing and discharging most of the cob back onto the ground. Early attempts to recover the cobs included adjusting the fingered vanes, on the harvester's chaffer, in order to open them wider and thus allow larger amounts of broken cob to fall through. However, this caused pieces of cob to lodge in the fingers and plug the chaffer. Accordingly, the vanes were lowered back to prevent plugging. Thus, only reduced volumes of cob residuals could be obtained for corn-and-cob-mix products.\nIn fact, in 1986, Chester McBroom was granted a patent on the Corn Cob Saver Sieve (U.S. Pat. No. 4,600,019). One drawback was the cobs could not be directly retrieved from the combine harvester. However, McBroom discovered, inter alia, that a cob-saving sieve structure need only have its inclined parallel vanes at an angle of 45°, and have the spacing between vanes 2″ apart, to allow shelled corn and harvested cob segments to pass therethrough. His goal was to provide a corn and cob segment mix.\nRetrieval of cobs directly, separately, and automatically from the combine was nevertheless undertaken by Larry Shrawder and disclosed in 1989 (U.S. Pat. No. 4,892,505) and then in 1993 (U.S. Pat. No. 5,256,106). Therein, combine harvest separator arcs were modified to a circular shape, from elliptical; and, the cross bar spacing on sieve slots was altered; and a cob conveyor was added. Again, the sieve slots were altered to a width of only 2-3 inches. His process and the special features of his combine design still limited harvesting to only cobs having lengths which did not exceed 2-3 inches.\nIn 1999, Vernon L. Flame, in U.S. Pat. No. 5,941,768 entitled Corn Cob Collecting Apparatus For Use With Combines, dealt with the problem of harvesting dry corn cobs. He found that the cobs could not be effectively separated from the stalks, husks, etc. through the use of air. Accordingly, he did not attempt to have cobs fall through the openings of a chaffer sieve. His solution was to use a wheeled apparatus connected to the rearward end of the combine in order to achieve separation of the cobs from a stream of crop material being discharged at the rearward end of the combine, where said end normally was designed to discharge the corn cobs, stalks, husks, etc. onto the ground.\nIn 2002, U.S. Pat. No. 6,358,141, by Stukenholtz, disclosed an on-board system for separating whole corn cobs by modification to the sieve and straw walker system. However, the system although allowing for the 2 inches between fingered vanes of the chaffer to allow the width of the cobs, offers no other sieve design modifications.\nIn recent years, as the demand and utilization of fossil fuels has skyrocketed, likewise has the demand for corn cobs as alternative renewable energy sources and as feedstock for environmentally friendly biodegradable packaging. Enormous volumes of corn cobs are needed as a principal feed stock for ethanol fuels. One company plans on harvesting cobs from plots having a minimum of 4000 acres as a component of their system that will eventually handle harvesting as much as 275,000 acres of cobs by as early as 2011. The cobs will be used, inter alia, by public utility companies as biomass to co-fire with coal, because the cob's burn value is competitive with that of wood.\nFor projects of this size and volume, cobs will be stored outside. Since broken cobs would more readily mold, and decay than whole cobs, the goal currently is to harvest the cobs as whole, or in lengths as long as possible. The industry therefore needs no cobs less than 3 inches in length, and as many whole cobs as possible. If the cobs were broken up, they would require storage in a dry environment, and much greater cost. One principal stumbling block to advancements in higher cob volume and higher pass throughs in sieve design was pointed out by Gordon in U.S. Pat. No. 7,011,579, issued in March 2006. Gordon pointed out that the proper opening between louvers is necessary, but if the opening is too large, straw and other material fall through contaminating the grain.\nAn improved chaffer and cleaning system for combine harvesters that would significantly enhance the recovery and harvesting of whole corn cobs, greater than 3″ in length, would fulfill a longfelt commercial need and provide a novel and unexpected technical advancement in the art."}
{"text": "In order to avoid heat losses through radiation, it is known to place a granular covering material on the top layer of a casting mold or ladle. The material used may be sawdust, glass, sodium silicate or Vermiculite. Recovered material such as aluminum oxide (Al2O3), calcium oxide (CaO) or aluminum drops may also be used.\nWhile the pouring of granular or powdered material onto the top layer of a casting ladle was initially carried out in a rather rough fashion by means of a hoisting mechanism and a chute, it was soon realized that it was necessary to find a more accurate solution. It is indeed desired to obtain accuracy in the metering of the desired quantity in order to obtain a predetermined layer thickness and also an even distribution of the bulk material over the whole surface to be covered.\nAn improved spreading device was suggested in EP 0 389 918, wherein the device comprises a fixed part from which a cover with scrapers is suspended, and a rotating part fitted below the scrapers. The rotating part comprises a series of slits that are uniformly distributed around a central shaft. The slits are radially widening in the direction away from the central shaft. Through the rotating movement of the rotating part and the levelling action of the fixed scrapers, the bulk material is uniformly distributed over the slits of the distribution bin. A rising and falling cover makes it possible to process varying volumes of bulk material.\nWhile this device achieves a homogeneous, protective and insulating covering layer on the top layer of a casting ladle containing molten steel it also has a number of disadvantages. Indeed, due to the central shaft, no material can be deposited directly to the center of the surface to be covered. This problem may however be partially solved by the use of deflectors for the bulk material.\nAnother disadvantage of the device is that after each operation, the distribution bin must again be filled with bulk material, which may cause some delay between two filling operations. Such a delay may e.g. cause unnecessary cooling of the molten steel or metal contained in the next ladle to be covered.\nFurthermore, it should also be noted that the construction of the device is rather cumbersome and necessitates a non-negligible amount of preventive and corrective maintenance work.\nAlthough, the device of EP 0 389 918 provides a homogeneous layer on the top layer of a casting ladle containing molten steel, there is still room for improvement."}
{"text": "1. Technical Field\nThe present disclosure relates to an ultrasonic sensor including a piezoelectric element for detecting ultrasonic wave and an acoustic matching member for conducting the ultrasonic wave received with an oscillation surface to the piezoelectric element.\n2. Description of Related Art\nPatent Document 1: JP-2009-20086A corresponding to US/20080072675A\nPatent Document 2: JP-2009-49776A corresponding to US/20090054784A\nSome ultrasonic sensors include a piezoelectric element and an acoustic matching member (see Patent Documents 1 and 2 for example). When the ultrasonic sensor is downsized, an oscillation surface (also called “a reception surface”) for receiving the ultrasonic wave and its sound pressure may become smaller. When the oscillation surface is made smaller, a receivable sound energy becomes smaller. Thus, it is difficult to maintain ultrasonic wave detection performance while the ultrasonic sensor is being downsized. A conceivable way for addressing the above difficulty may be the following. The acoustic matching member is formed so that an area of the oscillation surface is larger than an area of a region (also called “a contact surface”) contacting the piezoelectric element (see Patent Document 1 for example). In this configuration, the energy received with the oscillation surface is conducted to the contact surface, which is smaller in area than the oscillation surface. In this case, sound pressure concentration may occur and detention performance may improve. Thus, when the area of the contact surface is decreased so as to satisfy the above-described relationship between the oscillation surface and the contact surface, it is possible to maintain the ultrasonic wave detection performance while downsizing the ultrasonic sensor.\nIn the above related art, when the area of the oscillation surface is assumed to be constant, the ultrasonic wave detection performance improves with decreasing area of the contact surface. However, the piezoelectric element has a size limit from viewpoint of manufacturability. Thus, the above related art has a limit on the improvement of the detection performance. In addition, a decrease in the contact surface between the acoustic matching member and the piezoelectric element leads to a decrease in joint strength between the acoustic matching member and the piezoelectric element. With this regard, the above related art has a difficulty."}
{"text": "Various mechanisms are available to convert linear motion to rotational motion. Examples are ratchets, rack and pinion gears, and cranks. Almost all existing hydraulic manipulators use a slider crank mechanism to transform linear motion from a hydraulic actuator to rotary motion of a joint. There are two specific limitations to this approach. First, motion is generally limited to less than 120° of motion. Second, the transmission ratio (conversion of force to torque) is nonlinear and typically varies from about 2:1 to about 1:1 over the 120° workspace. The limitation of the transmission ratio requires that the actuators be oversized to give a specific torque requirement which adversely impacts the efficiency of a hydraulic system. Also, a slider-crank mechanism typically has a pivoting actuator to compensate for the circular motion of the crank center. Fluid powered actuators thus require some kind of flexible tubing or rotary union to pass fluids in and out of the actuator. These kinds of unions are much less reliable than a fixed porting and usually require considerable space. If, on the other hand, the actuator is not allowed to pivot, then an extra linkage is required to compensate for the circular crank motion. This additional linkage adds complexity and volume. What are needed therefore are more efficient and versatile mechanisms to convert linear motion to rotational motion."}
{"text": "U.S. patent application Ser. No. 612,208, filed Sept. 10, 1975, relates to a device which is attached to a vehicle hub for recording the distance covered by the vehicle. The device comprises a counting mechanism, a printing mechanism for stamp cards and feed means for advancing the counting mechanism wheels as a function of distance covered. The feed means is driven by an annular pendulum which is concentric with the vehicle wheel axle and is provided with mechanical coupling means, in the form of worm gears and screws on the pendulum axle, connected to the feed means for rotating the counting mechanism wheels.\nThe feed means for stepped feed of the counting mechanism wheels comprises a feed arm adapted to perform in response to the rotation of the worm gear in an arbitrary direction of rotation, a rocking motion between two stable end positions by action of a snap spring. The end positions are assumed by turning the feed arm by degrees to the dead point of the snap spring, whereafter the snap spring completes said rocking motion to the stable end position.\nThe counting mechanism wheels are provided with elevated numeral type faces for printing on the stamp cards. Consequently it is possible by inserting insert a thin, suitably bent metal sheet into the card slot and pressing it against an elevated numeral type, to press back the type. The counting mechanism can thus be manipulated to show a distance which is shorter than the actual distance covered by the vehicle."}
{"text": "1. Field of the Invention\nThe present invention is directed to an adjustable recoil apparatus for an automatic firearm.\n2. Description of the Prior Art\nIn an automatic firearm, such as an automatic pistol, a round within a firing chamber is detonated by impact from a firing pin. The explosion of the detonated charge produces a rapidly expanding volume of gas. The expanding gas propels a bullet from the round longitudinally through the barrel of the weapon. In addition, the expanding gas produces a force which propels the slide of the weapon located atop the pistol frame rearwardly relative to the frame. As the slide moves rearwardly an ejection finger or ramp ejects the spent shell from the firearm through the breech. In addition to effectuating ejection of the shell, the rearward movement of the slide cocks and feeds the next round to be fired into the slide. Successive rounds are typically provided from a magazine which, in the case of an automatic pistol, is normally located in the handle of the pistol frame.\nA recoil spring system is provided in an automatic firearm to return the slide forwardly along the pistol frame to position the next successive round in the firing chamber and to latch the slide relative to the frame. In conventional automatic firearms, such as automatic pistols, there is a single recoil spring located in the frame beneath the barrel. However, for some ammunition loads the recoil spring may provide too great a recoil force, and for other ammunition loads the force may be too light. If the recoil force is too great for the ammunition load the slide will not be fully propelled to the rear of the firearm frame. When this occurs the firearm will fail to eject the shell. On the other hand, if the recoil spring system is too light for the ammunition load there is an excessive amount of shock or \"kick\" within the firearm. This impairs the ability to accurately aim successive rounds during rapid firing.\nCurrently, in competition shooting a competitor will typically have at hand several alternative springs which can be interchangeably installed in an automatic firearm to provide the appropriate force for the ammunition load being used. However, in order to change recoil springs in the firearm the recoil system must be substantially disassembled. This results in a delay in shooting and also an inconvenience in locating an appropriate spring among a competitor's equipment."}
{"text": "The invention is applied to a device for non-invasive measurement of heart rate information, in particular to a heart rate monitor used in connection with exercise and sports.\nThe measurement of heart beat frequency is an interesting field of application in connection with exercise. On the basis of the heart beat frequency, i.e. the heart rate, it is possible to obtain information e.g. on a persons stress level, recovery and development of physical condition, and consequently the proportion of training exercises and rest can be monitored and planned better.\nThe heart rate is measured on a person\"\"s skin on the basis of an electrocardiographic (ECG) signal generated by a heart beat. Additional information on ECG is available in the following publication by Guyton, Arthur, C., Human Physiology and Mechanisms of Disease, third edition, W. B. Saunders Company, 1982, ISBN 4-7557-0072-8, Chapter 13: The Electrocardiogram, which is incorporated herein as reference. The electrocardiographic signal is an electromagnetic signal originating from a heart beat, which is detected on the body of a person to be measured. The signal is measured by means of electrodes, which are in contact with the body at least at two points. By a polarization vector, the electrode that is located closest to the heart often acts in practice as the actual measuring electrode, while the other electrode serves as ground potential, to which the voltage measured by the measuring electrode is compared as a function of time.\nThe heart rate monitor electrodes to be placed on the chest are arranged in a known manner in a belt-like structure, i.e. a so-called electrode belt. The electrode belt is thus a ring-shaped attachment means that goes round the whole chest and can be tightened round the human chest. A structure of this kind is shown in FIG. 1. The electrode belts are known to have structures that comprise an electronic unit in the middle of the belt, with an electrode on both sides. The electrodes measure the electric pulse transmitted by the heart and transmit the measurement results to the electronic unit through an interface connecting the electrode and the electronic unit. The components included in the electrode belt, such as the electronic unit and the electrodes, are generally coated with plastic or rubber in order to protect the components against moisture, for instance. Depending on the structure of the electrode belt, the electronic unit often also comprises means for transmitting an electric pulse as an analog burst to a receiver and display unit worn on the wrist, for instance. Alternatively, the electrode belt itself may comprise the means for storing and displaying the electric pulses.\nIn general, the electrode belts have a structure in which the rubber or plastic support structure covering the components of the electrode belt is relatively rigid. These electrode belts are, in general, poorly suited for long-term, continuous use, and they chafe the skin easily. The belt-like structure of the electrode belt also limits its optimal positioning substantially at the heart with persons having large quantities of muscular or other tissue in the chest area. Also slim adults and children have troubles in wearing the rigid electrode belt, because it does not bend sufficiently to follow the contours of a human body with a narrow chest. In some prior art solutions, the problem is approached such that the plastic support structure between the electronic unit and the electrode is pleated, whereby the electrode belt bends immediately outside the electronic unit. However, this solution only reduces rigidity in bending the electrode belt, because the electrode belt is still ring-shaped and attachable round the chest.\nThe prior art solution has a serious drawback: it is difficult to fit the rigid electrode belt optimally round the chest to achieve the best measurement result, in particular in long-term, continuous use, when the electrode belt also chafes the skin easily.\nThe object of the invention is to provide an improved electrode structure and heart rate measuring arrangement for measuring an electrical heart beat signal on a human body such that the above problems can be solved. This is achieved with the following electrode structure for measuring an ECG signal on the chest of a person. The electrode structure comprises a band-like component having an inner surface to be placed against the skin of the person\"\"s chest and an outer surface opposite thereto, and which electrode structure comprises a first electrode at a first end and a second electrode at a second end of the electrode structure, and the inner surface of the electrode structure is an adhesive surface for attaching the electrode structure to the skin of the person\"\"s chest, and the electrode structure is arranged to measure a potential difference between the first and the second electrodes caused by the ECG signal.\nThe invention also relates to a heart rate measuring arrangement for measuring the ECG signal on the skin of a person\"\"s chest. The heart rate measuring arrangement comprises an electrode structure placed on a person\"\"s chest and a wrist-worn receiver unit, the electrode structure comprising a band-like component having an inner surface against the skin of the person\"\"s chest and an outer surface, opposite thereto, and which electrode structure comprises a first electrode at a first end and a second electrode at a second end of the electrode structure, the inner surface of the electrode structure being an adhesive surface for attaching the electrode structure to the skin of the person\"\"s chest, and the electrode structure being arranged to measure a potential difference between the first and the second electrodes caused by the ECG signal, the electrode structure further comprising ECG processing means communicating with the electrodes for measuring the potential difference between the first and the second electrodes caused by the ECG signal and for producing heart rate information on the basis of the measured potential difference, and the electrode structure further comprising a transmitter for transmitting the heart rate information to the wrist-worn receiver which comprises a receiver for receiving the heart rate information transmitted from the electrode structure, the wrist-worn receiver further comprising a display for presenting the heart rate information.\nThe preferred embodiments of the invention are disclosed in the dependent claims.\nIn the solution of the invention, it is intended that the electrode structure is placed on the skin of the user\"\"s chest. In one embodiment, the band-like component of the electrode structure is of flexible, soft material that fits the skin closely, and as a consequence it is comfortable and inconspicuous to wear and does not chafe the skin, and hence it is well suited for long-term use in ECG measuring. In one embodiment the band-like component of the electrode belt is continuous, in which both electrodes and their attachment means are integrated. The band-like component is disposable and economical to manufacture. In terms of design, it is a plaster-like sticker, for instance.\nThe electrode structure has a first electrode at a first end and a second electrode at a second end. The first and the second electrodes of the electrode structure are electrically separated from one another in order to enable the measurement of the potential difference between the electrodes. For optimal measurement of the heart rate signal, the first and the second electrodes should be located sufficiently far apart from one another so as to detect an electric ECG signal generated by a heart beat. The electrodes are thus advantageously placed at the ends of the electrode structure. Naturally, there may be more than said two electrodes.\nAccording to a preferred embodiment, the inner surface of the electrode structure is an adhesive surface for attaching the electrode structure on the skin of the person\"\"s chest. An advantageous manner to implement the adhesive surface and the electrodes is the embodiment, in which the electrodes located at the ends of the band-like component of the electrode structure are made of electrically conductive adhesive. Hence, the adhesive attaches the electrode structure on the skin of the person\"\"s chest. In a second embodiment the first electrode and the second electrode of the electrode structure consist of an electrically conductive membrane at both ends of the electrode structure, at the electrode. On the inner surface of the membrane, which is placed against the person\"\"s skin there is an electrically conductive adhesive. The adhesive is preferably an electrically conductive glue. Further, a third manner to implement the electrodes and the adhesive surface is an embodiment, in which the first electrode and the second electrode of the electrode structure consist of an electrically conductive membrane at both ends of the electrode structure, at the electrode, and the electrodes at both ends of the electrode structure are narrower than the band-like component. Around the electrodes, on the outer edges of the band, on the inner surface thereof, there is an adhesive, with which the electrode structure is attached to the person\"\"s skin. In this case, the adhesive need not be electrically conductive. Because the inner surface of the electrode structure is that portion of the electrode structure which is against the person\"\"s skin, the electrodes are preferably located on the inner surface of the band. The electrode structure can also be designed such that the electrodes are partly or completely located on both the inner surface and the outer surface of the band. Further, the electrode structure can be designed such that the electrodes are located on the inner surface of the band, but they have interfaces also on the outer surface of the band.\nThe electrode structure also comprises an electronic unit communicating with the electrodes. The electronic unit is an electronic component attached to the band-like part of the electrode structure with one or more gripping means. The electrodes of the electrode structure communicate with the electronic unit. The electronic unit comprises ECG processing means, by which the potential difference caused by the ECG signal between the first and the second electrodes is measured, and an estimate for detected heart beat time instants is formed from the heart rate signals measured by the electrodes, and further the heart beat rate is calculated on the basis of the detected heart beat time instants. The electronic unit also comprises a transmitter for transmitting heart rate information to a wrist-worn receiver, which comprises a receiver for receiving the heart rate information transmitted from the electrode structure, and a display for presenting the heart rate information.\nThe wrist-worn receiver is located in a watch-like device that the user wears on his wrist, such as a heart rate monitor or a wrist computer. Transmission of information between the electrode structure and the heart rate monitor is thus carried out in known manners, for example through a connecting line, optically or electromagnetically. In the embodiment in question, the display for presenting the heart rate information is also preferably located in the wrist-worn receiver.\nThe electronic unit is preferably arranged in a casing which comprises one or more gripping means for attaching the electronic unit to the stap-like component of the electrode structure. The gripping means are most preferably located on that surface of the electronic unit casing which is against the person\"\"s skin. The preferable structures of the gripping means include attachment slots, a pivoted, clamping clip, or the like, that are in the central unit casing. The gripping means or the central unit casing comprises conductive connecting means, through which the ECG signal measured with the electrodes is applied from the electrodes to the electronic unit.\nThe invention also has an advantage that the electrode structure is inconspicuous, comfortable and well suited for long-term use as compared with the known solutions."}
{"text": "Conventional paperboard cartons for packaging milk, juice, and other beverage products are typically formed from a paperboard blank assembled in a rectangular open-ended configuration, which is then filled with the liquid contents and sealed with a gable or flat-folded top. The carton forming, filling, and sealing is done under sanitary conditions in lines of so-called \"form/fill/seal\" machinery. The industry has developed reclosable spouts for such cartons, typically in the form of plastic spout fitments which are sealed to the top end of the cartons. Examples of hinged-type spout fitments are shown in commonly-owned U.S. Pat. Nos. 4,705,197 and 4,770,325, both to Gordon and Kalberer. Other types of spout fitments include pull-open nozzles and screw caps. These reclosable spouts allow the cartons to be more easily opened and reclosed without leaking.\nThe plastic spout fitments are typically attached to the paperboard blanks by sealing a flange portion of the fitment to a heat-sealable coating or extrusion layer on the paperboard, by means of adhesives, heat sealing, or sealing with an ultrasonic horn. The attachment step requires an intermittent certain dwell time for aligning each fitment in position on a carton and applying the required adhesive, heat, or ultrasonic energy to the portion to be sealed. An example of indexing machinery for precise registration and sealing of fitments to cartons is shown in commonly-owned U.S. Pat. No. 4,846,915 of Keeler, Bombolevich and Sinocchi.\nAs reclosable spout fitments have come into widespread use, recent developments have been made in fitment sealing apparatus which synchronize the intermittent fitment sealing step with a continuous conveyor line used to supply cartons to form/fill/seal machinery. Examples of such fitment sealing apparatus are shown in copending, commonly-owned U.S. patent applications Ser. No. 07/594,536 of Keeler and Bombolevich, entitled \"Apparatus For Continuous Feeding And Synchronized Application of Fitments To Carton Blanks And Related Method\", and Ser. No. 07/571,811 of Keeler, entitled \"System For Continuous High Speed Application of Fitments To Carton Blanks And Related Method\", both of which are incorporated herein by reference. In these examples, the fitments are formed with a flat shape and are supplied to the fitment sealing station from a web or roll.\nIt is also desirable to apply three-dimensionally shaped spouts, e.g. plastic pull-open nozzles or screw caps, to the cartons at high speed in order to supply a continuous line of fitted cartons to a form/fill/seal machine. One proposal for such an apparatus is shown in U.S. Pat. No. 4,788,811 of Kawajiri, Honda and Furukawa. The Kawajiri et al. sealing apparatus includes a suction holder which retrieves a three-dimensional cap fitment from the end of a supply chute, moves the fitment axially inside a carton blank assembled in rectangular tube configuration, inserts the spout portion of the fitment laterally through a hole die-cut in the carton wall, and holds the flange portion of the fitment against the heat-sealable coating on the internal side of the carton while an ultrasonic horn is advanced on the external side. This apparatus, however, has the problem that the complex movements required of the cap holder limit the speed at which the machinery can be operated, and, further, has a significant risk of mechanical breakdown or misalignment of the fitment."}
{"text": "This invention relates to insecticidal butanoic acid esters.\nBelgian Pat. No. 801,946 discloses benzeneacetic acid esters as insecticides, and includes a commercially available compound, ##STR1##\nA substantial portion of the world's food supply is destroyed by pests and plant diseases. There is thus a continuing need for highly active insecticides."}
{"text": "A common design for a cube ice-making machine includes a vertical ice-forming mold. The mold has dividers that create individual pockets. When the pockets are sufficiently filled with ice, the control system for the machine switches into a harvest cycle. The ice cubes are released from the mold. The dividers may be sloped downward toward the open front so that the ice cubes slide out of the ice-forming mold under the influence of gravity, and into the ice collection bin.\nThe ice-making machine also includes a sump located beneath the ice-forming mold, a water distributor above the ice-forming mold, and a pump to transfer water from the sump up to the distributor. The water cascades down over the surface of the ice-forming mold. A part of the water freezes into the pockets and the rest runs off the surface of the ice-forming mold. A water curtain is placed adjacent to the ice-forming mold so that any splashing water is directed back into the sump. The bottom edge of the water curtain is bent to reach back under the ice-forming mold. This allows the front edge of the sump to be spaced behind the front of the ice-forming mold. With this design, the unfrozen water can return to the sump, but ice can fall straight down out of the ice-forming mold and into the collection bin.\nThe water curtain is typically suspended from pivots or hinges located near the top of the water curtain. The shape of the water curtain and location of the pivots are such that the center of gravity of the water curtain causes the sides of the water curtain to stay closed against the ice-forming mold frame while the machine is making ice. However, during the harvest cycle the water curtain can swing away as the ice is released from the ice-forming mold.\nA common technique for shutting down the ice-making machines when the bin is full is to place a sensor, such as a magnetic reed switch, near the water curtain, and put a magnet on the water curtain. The reed switch can then determine whether the water curtain is closed. This reed switch has two uses. First, when the water curtain closes, the machine can automatically switch back into an ice-making mode from a harvesting mode. Second, if ice has built up in the bin such that the slab of ice being harvested does not fall all of the way past the bottom edge of the water curtain, the water curtain will remain open, and the reed switch will not close. If the water curtain stays open for a sufficient amount of time, the ice machine shuts down and the “bin full” condition starts.\nDuring a freeze cycle, a thin bridge of ice forms over the dividers and between the individual cubes of ice. Most automatic ice-making machines allow for adjustment of the duration of the freeze cycle, which thus controls how thick this ice bridge becomes. A common control technique is to mount an ice thickness sensor so that as the ice bridge gets thicker, water running over the surface of it will contact a probe, directing the machine to automatically go into a harvest cycle. A thick ice bridge has the benefit that it helps in the harvest cycle, when water stops cascading over the front of the ice and the ice-forming mold is heated. A thick ice bridge allows the entire slab of interconnected ice cubes to be released at once. On the other hand, with a thin ice bridge, individual cubes have to each melt and drop out of their pockets, and adjoining cubes cannot help pull all of the ice out at once.\nWhile thicker ice bridges have some benefits, there are also some drawbacks. Because ice is an insulator, the efficiency of the freezing operation decreases as the ice bridge builds, since the heat is commonly transferred out of the back of the ice-forming mold by serpentine refrigerant coils forming the evaporator section of a refrigeration system. Most importantly, many end users do not want thick ice bridges, because the slabs of ice cubes do not break into individual cubes as easily, and chunks of ice cubes frozen together are hard to dispense, scoop or fit into a cup.\nAccurate control of ice bridging can be effective in providing properly sized ice cubes for the user. Care must be taken, however, to avoid the inadvertent freezing together of the ice cubes once they are delivered to the ice bin. Water going down the face of the ice-forming mold may fall into the ice bin rather than being directed by the bottom of the water curtain back into the sump. The leakage of water into the ice bin can cause the ice cubes in the ice bin to freeze together, or create wet ice in the ice bin. Also, for those machines that do not add water during the freeze cycle, and go into harvest when the water level drops to a predetermined point, the loss of water from the water recirculation system will result in less ice being made in each cycle. Accordingly, a need existed for an improved water curtain and ice machine that minimizes the leakage of water into the ice bin of an ice machine."}
{"text": "PTL 1 describes an example of the above-described cutting insert according to the related art. The cutting insert disclosed in PTL 1 includes a substrate having a seat at a corner thereof, the seat being recessed from the top surface of the substrate. A superhard sintered body is bonded to the seat, and a cutting edge, a rake face, and a chip breaker are formed in the superhard sintered body.\nIn the cutting insert according to PTL 1, a relatively large negative rake face is formed along edges of the top surface of the substrate and the superhard sintered body. A flat land, which functions as a rake face, is formed in the superhard sintered body by removing a corner of the superhard sintered body along a ridge line that function as a cutting edge at a position such that the negative rake face somewhat remains. A chip breaker is formed by cutting the superhard sintered body obliquely upward from an end of the land.\nThe chip breaker includes a projection that projects from the position of the negative rake face toward a tip end of an arc edge such that the amount of projection is at a maximum at the center thereof and a breaker wall that extends obliquely upward from the rake face toward a top edge of the projection. The breaker wall includes two surfaces that are connected to each other at an angle."}
{"text": "Existing work surfaces made from concrete, such as a kitchen work surface, are manufactured on site in the kitchen. The concrete is poured into a bespoke trough mold, leveled, and allowed to cure."}
{"text": "(1) Field of the Invention\nThe invention relates to three-dimensional imaging. More specifically, the invention relates to capture and distribution of three-dimensional digital images.\n(2) Background\nVarious non-contact methods for creating three-dimensional models of three-dimensional objects exist. There are four broad categories of non-contact three-dimensional digitizers. The first category is referred to as silhouette digitizers because the imaging device repeatedly takes the silhouette of the object as the object is rotated before the imaging device or the imaging device is rotated about the object. This type of digitizer is relatively ineffective at dealing with concavities in a three-dimensional object because the silhouette is unchanged by the concavity.\nThe second category is timing digitizers. Timing digitizers use a signal source such as a radar source. By determining the amount of time required for the signal to bounce back from the different points on the target object, surface features of the object are revealed. However, such digitizing methods are extremely expensive to implement.\nA third category is projected pattern digitizers, in which a pattern of some sort is projected onto the object to be digitized and the dimensions of the object are determined from the interaction of the pattern with the object. Projected pattern digitizers fall into three main subcategories. The subcategories include contour digitizers which use spatial distortion from a projected pattern of contour lines to determine surface features of a three-dimensional object. A next subcategory is interference projected pattern digitizers, which use two sources and then based on the localized interference pattern of the two sources, determine the surface features of the three-dimensional object to be digitized. A third subcategory is referred to as color projected pattern digitizers because this category uses a projected color pattern and resulting color gradients to determine relevant information about the object to be digitized.\nA final broad category is stereoscopic digitizers which employ multiple cameras to capture images of the object from different angles. From the picture, such systems perform feature identification. Then a correlation between the features in the different pictures is established to yield three-dimensional data.\nCommon to all of these various methods are high cost and that each typically mandates a highly constrained environment often requiring a lab-like setting in which to perform the imaging in order to attain necessary alignments. This cost and difficulty of use render such systems impractical for use for the average user."}
{"text": "The present invention relates generally to apparatus for the application of adhesives or solvents to the surfaces of thermoplastic materials to be bonded together, one to the other, and more particularly to apparatus for the spray application of adhesives, either as film forming materials held in a solvent solution or as solvents per se to surfaces of thermoplastic materials, such as polystyrene in either sheet or foamed form, with the spray application being controlled within a spray enclosure, and being controllably applied to surfaces of the workpieces to be bonded. The present apparatus is particularly adapted for use in connection with solvent bonding of polystyrene, including foamed polystyrene, and particularly wherein solvents such as trichloroethylene or other more dangerous solvents may be employed, and wherein the atmosphere from the spray enclosure is controllably monitored so as to avoid contamination of the work area, and also to avoid pollution of the environment.\nSolvent bonding of workpieces such as surfaces of typical thermoplastic articles including such articles in sheet form, complex shaped articles such as vacuum-formed articles or molded articles, and including solid as well as foamed thermoplastics has been known and highly utilized in the past. Typically, the surfaces to be bonded are treated with an application of adhesive in a solvent solution or a solvent per se so that the surfaces of the workpiece portions may be bonded together in face-to-face relationship. The means of solvent and adhesive application including spray application, brush application, dipping and the like. Typically too, and in solvent application operations, solvent has been applied without regard to the quantity of solvent employed, the manner of disposing of the solvent, or the manner in which the environment is protected from contamination due to the inevitable vaporization or volatilization of the solvent.\nSolvent bonding of foamed articles has generally been regarded as an extremely difficult or, in some instances, an impossible operation. However, in accordance with the aspects of the present invention, it has been found possible to bond foamed polystyrene workpieces together without adversely affecting the quality of the product and with a firm bond created. Such operations are made possible in accordance with the various aspects of the present invention."}
{"text": "A computer or an electronic system usually requires storage devices for storing the data and information that it processes and needs to support its operations. These storage devices may be internally mounted within the housing or external to the computer or electronic system.\nAs computer and electronic system components become smaller, the enclosures that house these systems also become smaller. The advantage of reducing the size of such systems is for ease of transportation, installation, reduced total power consumption, and reduced office space consumption.\nAny given computer system is designed to fit any variety of customer applications. Some of these applications require relatively little disk storage and others a great amount. The object of any computer design is to fill as many needs as reasonably possible without making the system excessively large or expensive. For many applications, this results in the need for varying amounts of external storage. Directly connected external storage devices are usually located adjacent to the computer or electronic system, but even short cables negatively impact the reliability and performance of such external storage devices.\nTo properly protect the environment from contamination from noise contained within the cables connected to the external storage devices, and to protect that data within the cable from possible corruption from outside noise, the cable run must be shielded in some way. However, such shielding significantly contributes to the difficulty in designing a reliable and cost efficient electrical data connection.\nOne of the advantages, however, in using external storage devices is that there is no practical limit in the number of storage devices that may be connected to the computer or electronic system. Therefore, as long as there are unused input/output (\"I/O\") ports, storage devices may be added to increase the system's storage capacity.\nComputer housings have shielding associated with them. Therefore, when storage devices are internally mounted, interference and noise internal to the housing that they generate will not radiate to the external environment. However, there are a number of problems associated with internally mounting storage devices. A few of these problems are the limitation on the number of storage devices that can be mounted to a storage assembly because of the space constraints, the existence of lengthy electrical stubs associated with a common bus, difficulty in providing clean power to the storage devices, the ability to easily expand a system's storage capacity, the difficulty in cooling the storage devices, and the ability to isolate the storage devices from external sources of physical shock and vibration.\nConventional storage assemblies generally do not employ methods which permit more than a few storage devices to be internally mounted on a storage assembly. For example, some conventional storage assemblies only permit one storage device to be connected to them. In many cases it is desirable to provide a storage assembly that accommodates the connection of a number of storage devices.\nInternally mounted storage devices, which are conventionally connected to a computer or an electronic system, often have lengthy stubs associated with the connection of storage devices to a common bus. Because of the transmission line effects, the introduction of lengthy stubs limit the maximum configuration and reliability of the common bus that connects to all of the storage devices. It is desirable to have a storage assembly that is not plagued with the stub problem.\nModern storage devices require high current and high speed changes in current. These operational power needs for internally mounted storage devices are not satisfied with conventional techniques for providing power because such operational power needs of modern storage devices exceed the capability of conventional power distribution methods. Thus, it is desirable to have a storage assembly that can satisfy the required operational power needs of modern storage devices.\nThere also has been difficulty is expanding the storage capacity of computer or electronic systems beyond that which is provided by the internally mounted storage devices. It fact, many systems with internally mounted storage devices are not expandable at all.\nThe expansion problem also is manifest in individual storage assemblies being part of a larger storage assembly with all of the individual assemblies being connected by a common bus. Such a configuration cannot be readily expanded because of the difficulty in making bus and other connections without adversely affecting the system.\nThe present invention provides an storage assembly that overcomes these and other problems as will be set forth in detail in the remainder of the specification and shown in the drawings."}
{"text": "Recently, input devices having a display pattern disposed in an outer area located outside an input area for the purpose of improvement in design of the input devices have been known (for example, see Patent Literature 1). Such an input device includes a substrate and a protective member bonded to the substrate with a bonding member interposed therebetween. The input device includes coloring members for plural colors covered with the bonding member between the substrate and the protective member, whereby a display pattern is formed in the outer area."}
{"text": "1) Field of the Invention\nThis invention relates to a tethered buoy housing and deployment assembly adapted to lift and rotate a buoy from a stowed location in an undersea vessel to a position wherein the buoy can be released into an underwater flow stream and thereafter retrieved and returned to the housing.\n2) Description of the Prior Art\nThe United States Navy has developed an antenna assemblies for submarines, in which the assemblies are adapted to improve communications at maneuvering speeds and depths. One such system is the Recoverable Tethered Optical Fiber (RTOF) buoy, which is deployed to the surface and recovered by an attached tether. Current RTOF buoy system is designed to fit within a relatively large working volume. As such there is a need for a buoy system which can operate in a relatively smaller space."}
{"text": "With the rapid development of the Internet, computer users all over the world are becoming increasingly more exposed to writings that are penned in non-native languages. Many users are entirely unfamiliar with non-native languages. Even for a user who has some training in a non-native language, it is often difficult for that user to read and comprehend the non-native language.\nConsider the plight of a Chinese user who accesses web pages or other electronic documents written in English. The Chinese user may have had some formal training in English during school, but such training is often insufficient to enable them to fully read and comprehend certain words, phrases, or sentences written in English. The Chinese-English situation is used as but one example to illustrate the point. This problem persists across other language boundaries.\nAccordingly, this invention arose out of concerns associated with providing machine-aided reading systems and methods that help computer users read and comprehend electronic writings that are presented in non-native languages."}
{"text": "1. Technical Field\nAspects described herein relate to a communication device, a communication system, a communication method, and a computer readable recording medium configured to configure data transmission conditions in accordance with a network state and to perform communication under the configured transmission conditions.\n2. Description of the Related Art\nCommunication devices capable of adjusting data transmission conditions in accordance with the network state and performing communication under the adjusted data transmission conditions are known. In the related art, for example, a communication device that adjusts a data transmission period in accordance with a communication load on a network has been proposed. The communication device is capable of reducing a data transmission period when the network bandwidth usage rate is low.\nIn some examples, a relay device that manages relay and transfer of data is located in a network. The relay device temporarily stores data received from a communication device. The stored data is relayed and transferred to another communication device at predetermined intervals. In this case, even if a communication device reduces the transmission period of the data, the relay device relays and transfers data at constant or consistent intervals. Thus, the communication between the communication devices may be delayed because the data transmission speed and timing depends on the constant/consistent transmission intervals of the relay device."}
{"text": "1. Field of the Invention\nThe present invention relates to a control apparatus of an engine having a variable valve lift apparatus and a method thereof. More particularly, the present invention relates to a control apparatus of an engine having a variable valve lift apparatus and a method for constantly maintaining pressure of a cylinder by driving an intake valve or an exhaust valve in a low lift mode when a cylinder is in an idle state.\n2. Description of Related Art\nIn general, an internal combustion engine generates power by taking in fuel and air and burning it in a combustion chamber. Intake valves are operated by a camshaft in order to take in air, and air is drawn into the combustion chamber while the intake valves are open. In addition, exhaust valves are operated by the camshaft, and exhaust gas is discharged from the combustion chamber while the exhaust valves are open.\nAn optimal operation of the intake valves and the exhaust valves depends on the rotation speed of the engine. That is, an optimal lift or an optimal opening/closing timing of the valves depends on the rotation speed of the engine. In order to achieve such optimal valve operation depending on the rotation speed of the engine, research has been undertaken for a variable valve lift (VVL) apparatus that includes a plurality of cams for driving the valves, or varies the lift of the valves depending on the number of rotations per minute of the engine.\nIn a similar concept to the VVL apparatus, there is a CDA device. In general, cylinder deactivation (hereafter referred to as CDA) is a technology of deactivating some cylinders when braking or driving at a predetermined speed. During the operation of the CDA device, fuel supply to the deactivated cylinders is stopped and the operation of the intake/exhaust valves is stopped.\nWhen a cylinder is deactivated, the intake valve and the exhaust valve are both closed, and an explosive gas or intake air is trapped in the cylinder. When the explosive gas or intake air is trapped in the cylinder, as shown in FIG. 1 (refer to arrow), there is a problem that an impact or a vibration is generated by a rapid increment of pressure in the cylinder at an initial state of cylinder deactivation.\nAs time passes, the pressure leaks outside of the cylinder, and thus negative pressure is formed in the cylinder. Therefore, oil flows into the cylinder.\nFurther, since the pressure in the cylinder is recovered to an original state when the cylinder is activated, there is a problem that fuel consumption is deteriorated.\nThe information disclosed in this Background of the Invention section is only for enhancement of understanding of the general background of the invention and should not be taken as an acknowledgement or any form of suggestion that this information forms the related art already known to a person skilled in the art."}
{"text": "The expression of recombinant polypeptides in host cells is a standard technology widely used in biotechnology and the pharmaceutical industry. Especially microbial hosts, like for example E. coli, are commonly used, since relatively simple expression systems and cell culture conditions are available for these host cells. In general the cultivation process is therefore comparably economic.\nIt is, however, often difficult to obtain a polypeptide of interest in soluble and active form when it is expressed in microbial cells. Often, expression of a recombinant polypeptide leads to the production of poorly soluble intracellular aggregates of the polypeptide in denatured form, the so-called inclusion bodies [Baneyx, F. and Mujacic, M. (2004) Nat. Biotechnol. 22, 1399-1408 and Sorensen, H. P. and Mortensen, K. K. (2005) Microb. Cell Fact. 4, 1], also referred to as classical inclusion bodies.\nClassical inclusion bodies are generally easy to isolate, typically by centrifugation at moderate speed. To recover the active, i.e. correctly folded, polypeptide from the inclusion bodies, the inclusion bodies have to be solubilized and the protein renaturated after isolation. Several patent applications and patents deal with the aspect of solubilizing the inclusion bodies and renaturing the proteins obtained from inclusion bodies. For example EP0512097, EP0364926, EP0219874, WO01/87925, Rudolph 1996, Rudolph 1990, Marston 1986 and Dietrich 2003 describe general techniques relating to the solubilization and renaturing of denatured proteins. For example, EP0219874 discloses generic methods for refolding of recombinant proteins from E. coli inclusion bodies. For the solubilisation the chaotropic agents GuHCl and arginine were used at high pH. EP0219874 describes the formation of disulfide bridges under redox conditions provided by GSH/GSSG.\nDespite the fact that numerous processes for the isolation and solubilization of inclusion bodies are known, the results are not always satisfactory. One major problem is that the structure of the inclusion bodies can vary. It is known that the formation and the structure of the inclusion bodies can be influenced by parameters of the cell culture process, including for example media composition, growth temperature and production rate. WO2004/015124 describes the formation of “non-classical” inclusion bodies by modulating culture conditions.\nObtaining recombinant proteins from inclusion bodies, in particular in active form and sufficient amounts, can be problematic. Sometimes the structure of the inclusion bodies is too “soft” which leads to the situation that the isolation of the inclusion bodies by centrifugation is difficult. On the other hand it is also possible that the inclusion bodies are too compact. This results in inclusion bodies which cannot be solubilized even under rough conditions.\nTo overcome these disadvantages the present invention provides new cell culture processes which lead to large amounts of properly folded protein in inclusion bodies that can easily be isolated and solubilised, resulting in increased yield of the recombinant protein."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of televideo communications and, more specifically, to two-way video teleconferencing equipment used in the office environment.\n2. Description of the Background\nA video teleconferencing unit has been proposed for use in the office environment in connection with a personal computer (\"PC\") to enable one party at a first location to visually and audibly communicate with a second party at a second location. The sound is detected by a microphone and transmitted and reproduced by a loud speaker. The visual image is picked up by a lens of a video camera and reproduced on the monitor of the PC.\nIt is becoming more desirable to use televideo communication equipment in the office environment, as it is an efficient mode of communication with innovative features such as simultaneous visual and audio transmission. The communication systems currently available, however, typically consist of several individual components, such as a camera housing, a remote microphone, and separate speakers, thereby rendering the system unduly cumbersome.\nFurthermore, in current teleconferencing systems, the lens assembly is fixedly attached to the main housing of the camera and transmitting unit. Therefore, in order to adjust the angle or position of the lens, the entire main housing must be moved or propped up to achieve varying angles. Such movement is undesirable as it leads to possible damage to the teleconferencing unit from excessive movement and it decreases the adaptability of the unit."}
{"text": "1. Technical Field\nThis disclosure relates to image capturing, and more specifically, to the compression of image data in an image sensor chip, prior to transmission to a digital signal processor (DSP).\n2. Description of the Related Arts\nThe advancement of digital video and image encoding has led to increasingly sophisticated image capture techniques at increasingly high resolutions and frame rates. For instance, common media formats, such as Blu-Ray discs, internet video, and cable/satellite television, are able to display content at a 1080P resolution (1920×1080 pixels progressively scanned) at 60 frames per second (“fps”). Certain displays are able to display resolutions of 2560×2048 pixels or 3260×1830 pixels or higher at frame rates of 120 frames per second or higher. As encoding and display technology advances, frame rates and resolutions will increase accordingly.\nThe capture of digital images by an image capture device (hereinafter “camera”) is performed by an image sensor. Many types of image sensors are commonly used in cameras and other image-capturing devices, such as charge-coupled devices (CCDs) and complementary metal-oxide-semiconductors (CMOSs). Image sensors convert light, such as light entering the aperture of a camera through a camera lens, into image information. In this way, a camera can “capture” objects before it by converting the light reflected from the objects and passing through the camera lens into an image.\nImage data from image sensors must often be processed into a particular image format prior to the image being rendered by a display. Image data is typically processed by a DSP, which is often located off-chip from the image sensor. Image sensors and DSPs are typically connected by buses coupled to the image sensors' pins and the DSPs' pins. As image resolution and frame rate increase, the amount of image data produced by the image sensor increases, along with the amount of power required to accommodate the transfer of image data from the image sensor to the DSP. In some circumstances, the amount of power required to accommodate such a transfer of image data rapidly depletes the battery life of a camera, or exceeds the amount of power available to the camera. In addition, increasing the power consumption in transferring image data correlatively increases the noise (such as the electromagnetic noise) impacting the image data.\nThe limited ability to transfer image data between an image sensor and a DSP thus creates a bottleneck in the image capture process. One solution is to increase the number of pins on the image sensor and DSP to increase the bus bandwidth between the image sensor and DSP. Such a solution, however, may not be physically possible due to the limited real estate of many image sensor chips and DSPs. Alternatively, images or video may be captured at a lower resolution and at a lower frame rate to reduce the power-limiting bandwidth requirements between the image sensor and the DSP. However, this solution results in a lower quality image or video. Thus, camera makers are increasingly required to balance the competing requirements of image and video quality with power consumption and chip real estate."}
{"text": "This invention relates to galvanic battery cells and more particularly to sealants applied to such cells in their manufacture to prevent the escape of electrolyte from the cells.\nThe problem of electrolyte leakage from battery cells is well known. In attempts to overcome this problem, various mechanical solutions have been proposed such as seal rings and gaskets (e.g. U.S. Pat. No. 3,068,312, issued to Daley et al shows a seal ring and U.S. Pat. No. 3,068,313, also issued to Daley et al shows a seal gasket).\nTypically, such seal rings have been coated with a dielectric material. As those skilled in the art are aware, the coating should be soft enough to deform so as to fill the imperfections or gaps between the cell container and the sealing ring or gasket, thereby shutting off all significant leakage of electrolyte.\nPrior to the present invention, a common battery sealant consisted of a mixture of three parts (by volume) of naphtha solvent and one part (by volume) of Pioneer asphalt paint. Such a sealant has been commercially marketed by the Pioneer Division of the Witco Chemical Company of Perth Amboy, New Jersey. The Pioneer asphalt paint has been marketed under the name Pioneer 609 paint and is in accordance with Federal Specification MIL-C-450B Type II; TT-C-494 Type II.\nThe naphtha in the sealant provided by Pioneer meets the specifications of Federal Specification TT-N-95B Type I.\nAlthough the Pioneer 609/naphtha solution has functioned adequately in battery cells to prevent the leakage of electrolyte, its application has involved problems because of the flammability of the solution.\nIt is therefore an object of this invention to provide a battery cell sealant which prevents electrolyte leakage but which is less flammable than sealants containing naphtha.\nIt is a further object of this invention to provide a battery sealant of low flammability which has unexpectedly improved sealing characteristics."}
{"text": "The present invention relates to a new composition for transdermal delivery of topically applied pharmaceutical preparations. The system comprises the use of the pharmaceutical agent and acidified nitrite contained within a delivery system to allow passage of both the specific pharmaceutical agent and nitric oxide to the skin.\nThe penetration of substances through the skin is important from both toxicological and therapeutic viewpoints. Passive delivery of most compounds across different epithelia is limited due to the excellent barrier properties afforded by these epithelia. The stratum corneum is the principal barrier to penetration of most chemicals. Conventional topical delivery systems are therefore restricted to either substances for local effects or to highly potent, small, lipophilic substances for systemic effects. It is also difficult to deliver ionic and high-molecular-weight drugs in therapeutically sufficient amounts by conventional systems.\nBy way of example, many medical and surgical procedures require topical anaesthesia. The use of local anaesthetics requires an agent possessing the following general properties. It should not be irritating to the tissue to which it is applied, nor should it cause any permanent damage to nerve structure. Its systemic toxicity should be low because it is eventually absorbed from its site of administration. It is usually important that the time required for the onset of anaesthesia should be as short as possible. Furthermore, the action must last long enough to allow time for the contemplated medical or surgical intervention, yet not so long as to entail an extended period of recovery (J. Murdoch Ritchie & N. M. Greene, Local Anaesthetics in Goodman & Gilman's: The Pharmacological Basis of Therapeutics, pages 311-331, McGraw-Hill Inc. (1992)).\nLocal anaesthetics are rapidly absorbed into the circulation following topical administration to mucous membranes or denuded skin. It is extremely useful in achieving loss of sensation in a subject without the loss of consciousness, or the impairment of central control of vital functions. Typical uses, include minor invasive procedures such as venepuncture, e.g. for the collection of blood for diagnostic purposes from a patient, for the administration of therapeutic agents, whole blood or blood plasma to a patient, or prior to the administration of a general anaesthetic to a patient. However, it is a common feature among patients that the pain of injection can cause discomfort and in certain cases a patient, in particular juvenile subjects, can experience acute anxiety or panic brought on by the sight of a needle or of the injection itself. Such panic attacks can be characterised by fainting, vomiting or other related symptoms. Whether the adverse reaction is pain or a panic attack, the problem leads to poor patient compliance with advisable medical procedures. There exists a need therefore for improved local anaesthetic compositions that can overcome these problems.\nIntact, healthy human skin presents an excellent natural barrier to the external environment and restricts the passive diffusion of pharmaceuticals. Local anaesthetics do not readily penetrate intact skin (McCafferty et al., Br. J. Anaesth. 60, pages 64-69 (1988)).\nThe insertion of a needle through the skin, for procedures such as phlebotomy or vaccination, is painful and may induce great fear and anxiety especially in children and the elderly. Painful experiences lead to reduced compliance, with heightened anticipatory anxiety and fear. The introduction of topically-applied cutaneous anaesthetic preparations such as EMLA™ (Eutectic Mixture of Local Anaesthetics), cream [Astra Pharmaceuticals Ltd.] (Arts et al., Pediatrics 93, pages 797-801 (1994)) and more recently Ametop Gel™ [Smith & Nephew Healthcare Ltd.] (Freeman et al., Paediatr. Anaesth. 3, pages 129-138 (1993)), represented a definite advance in clinical practice and are contributing to breaking the cycle of “needle phobia.”\nStudies of the effects of these preparations have produced variable results (Moiodecka et al., Br. J. Anaesth. 72, pages 174-176 (1994); Lawson et al., Br. J. Anaesth. 75, pages 282-285 (1995)). However, relatively slow onset times (EMLA™ 60-90 minutes; Ametop Gel™ 30-45 minutes) remain a deterrent to widespread clinical and patient acceptance with the need to organize clinic, ward and operating theatre routines accordingly. These methods are of no benefit in acute situations. Additionally, even following the manufacturer's recommendations for dosage and administration, potential exists for improvement in the degree of anaesthesia afforded by these treatments.\nA percutaneous local anaesthetic with a more rapid onset time and increased potency would be helpful in organisational terms for emergency cases, community medicine and for an increasing number of paediatric medical and surgical day cases. Shortening of the period of anticipatory anxiety while achieving the maximal desensitising of the skin would clearly be clinically advantageous.\nNitric oxide [NO] is a potent vasodilator synthesized and released by vascular endothelial cells and plays an important role in regulating vascular local resistance and blood flow (Palmer et al., Nature 327, pages 524-6 (1987)). In mammalian cells, NO is produced along with L-citrulline by the enzymatic oxidation of L-arginine. Nitric oxide is also involved in the inhibition of both platelet and leukocyte aggregation and adhesion, the inhibition of cell proliferation, the scavenging of superoxide radicals and the modulation of endothelial layer permeability. Nitric oxide also has been shown to possess anti-microbial properties, reviewed by F. C. Fang (1997) (J. Clin. Invest. 99 (12):2818-2825 (1997)).\nA potential therapeutic utility of the anti-microbial properties of NO is described in WO 95/22335. A pharmaceutical composition comprising nitrite in an inert carrier cream or ointment and salicylic acid was used to show killing of cultures containing E. coli and C. albicans. This activity was further tested against patients with fungal infection of the feet (“Athlete's Foot” or Tidea pedis) and showed that the condition was amenable to treatment with the acidified nitrite composition. However, the composition of nitrite and organic acid caused erythema (redness) of the skin.\nIn addition to internal cell-mediated production, NO is also continually released externally from the surface of the skin by a mechanism, which appears to be independent of NO synthase enzyme. Nitrate excreted in sweat is reduced to nitrite by an unknown mechanism, which may involve nitrite reductase enzymes, which are expressed by skin commensal bacteria. Alternatively mammalian nitrite reductase enzymes may be present in the skin which could reduce nitrite rapidly to NO on the skin surface (Weller et al., J. Invest. Dermatol. 107, pages 327-331 (1996)).\nThe production of NO from nitrite is believed to be through the following mechanism:NO2−+H+<=>HNO2  [1]2HNO2<=>N2O3+H2O  [2]N2O3<=>NO+NO2  [3]\nTopical application of a sodium nitrite/ascorbic acid NO-generating system causes significant increases in skin blood flow in patients with Raynaud's disease and in normal healthy subjects without causing local irritation (Tucker et al., Lancet 354(9191): 1670-5 (1999); Harwick et al., Clinical Science 100, pages 395-400 (2001)). The reaction can be terminated within a few seconds by gentle wiping of the skin with a tissue."}
{"text": "This invention relates generally to bayonet joints of the type wherein a male portion is retained within a female portion and in particular, to a rigid axial retainer for use in retaining the male portion within the female portion.\nIt is well known that most modern handles for brooms, similar cleaning equipment and the like are made of lightweight metallic tubing. Such handles are often covered with plastic and are widely distributed.\nThese handles commonly include a plurality of sections which can be combined with each other in order to vary the length of the handle and permit the length of the handle to be changed for each specific need or use to which the handle is to be put. The ability to vary the length of the handle also facilitates packing, storing and transportation of both the handle and the equipment with which the handle is to be used.\nThe system of connecting sections of handles together and of connecting the handles to an attachment such as a broom or other cleaning device is normally based on so-called \"bayonet joints\". As used herein, the term \"bayonet joints\" refers to joints which include a \"male\" length or member with a suitably tapered end which is adapted to be inserted into a \"female\" length or member. In order to connect a series of male and female lengths, each length has one tapered end which functions as a male end and one untapered end which functions as a female end. The lengths can then be connected one after the other by alternately fitting male ends into female ends to achieve any desired overall length.\nThere are many systems for securing the male and female lengths so that the bayonet joints do not become undone during use. In general, means are provided for securing the lengths in a longitudinal direction, hereinafter referred to as the \"axial\" direction.\nIn one embodiment, the taper of the male length is provided with a thread which can be screwed into a corresponding internal thread in the female length. The threads are produced by rolling the lengths directly on a thread-forming plate. It is necessary for the threads to be tight in order to impart rigidity to the joint. The difficulty with constructions of this embodiment is that the firmness of the thread spiral decreases after repeated threading and unthreading of the male length in and out of the female length. The decreased firmness that develops tends to pull the male length out of the female length and cause the male length to act against the female length with a wedge action. This in turn causes the edges to be deformed to such an extent that it is no longer possible to tighten the threads properly, and the joint loses rigidity.\nIn another similar embodiment, the threads of the male length and the female length are provided on separate attachments which are then separately secured to each of the male and female lengths. However, this alternative necessarily increases the weight and the cost of the system.\nIn still another embodiment, the male length is provided with a circumferential hole. A rounded pin slides radially through the hole against the action of an opposing spring provided within the male length. During engagement of the male length with the female length, the pin retracts within the taper of the male part. Following engagement, the pin is pushed by the spring and extends through the circumferential hole in the male part and a corresponding hole arranged radially on the female part, thereby axially locking the joint. This system does not, however, avoid the drawbacks associated with radial or transverse play of the joint.\nIn order to eliminate the widespread drawbacks associated with lack of rigidity in units connected by the devices provided in accordance with the known art summarized above, the device disclosed and claimed in European patent application No. 84 83 0295.6, was created. This device has a cylindrical expansion ring on the male part. The ring expands when introduced into the female part by axially striking against a stop provided in front of the male piece under the action of a threaded conical frustum, which then penetrates into the ring itself. The expansion ring is further provided with a longitudinal slit in order to favor homogeneous deformation of the ring. This device provides excellent axial holding characteristics to bayonet joints in which it is used and does so without permitting any transverse or radial play. Consequently, optimum rigidity is provided. The corresponding unit using this device in a bayonet joint is, however, generally too expensive to permit use in brooms, cleaning equipment and the like."}
{"text": "I. Field\nThe present disclosure generally relates to network devices. More particularly, the disclosure relates to multi-homing network devices.\nII. Description of Related Art\nIn recent years, the number of wireless handsets in operation has increased dramatically. With the increase in demand for these mobile devices, the manufacturers are building these devices to include numerous data services. This convergence of data services can make wireless devices powerful resources for data networking. However, due to the increase in data services provided by wireless devices, the number of nefarious attacks on wireless devices has also increased. Hence, there is an increasing concern to protect these devices from malicious attacks.\nWireless telephones have become multi-homing devices that include many data interfaces through which the wireless telephones can accept and send data. At any time, within a particular wireless telephone, more than one of these data interfaces can be open for direct data transfer with the Internet, or another data network. The data protocol stack in the mobile device is mostly transparent to the multiple data interfaces of the device. Further, the data protocol stack can accept data from any of the data interfaces as long as the protocol address of the incoming data matches the protocol address of the phone. As such, the wireless telephone is open and vulnerable to many attacks from the Internet and other data networks.\nFor example, when a packet is received on a data interface at a multi-homing device, the packet can be routed to an appropriate socket, or application. In general, for a socket that is connected using transmission control protocol (TCP), a packet is routed to the socket based on four (4) tuples, e.g., source address (src_addr), source port (src_port), destination address (dst_addr), and destination port (dst_port). For a socket that is connected using user datagram protocol (UDP), a packet is routed to the socket based on two tuples, e.g., destination address (dst_addr) and destination port (dst_port). For other protocols, such as Internet control message protocol (ICMP) or non-Internet Protocol (IP) based protocols, other fields in the network and transport layer headers can be used.\nUnfortunately, in a multi-homing device the parameters described above may not be able to uniquely identify a data interface for several reasons. For example, the data interfaces within a multi-homing device may be assigned duplicate private addresses. Also, multiple applications within the multi-homing device may try to access the same service using different network data interfaces. In such scenarios, the applications may bind to the same service access point (SAP), e.g., the same port number in case of UDP or TCP. As such, it may not be possible to correctly route a packet to the appropriate destination application.\nA multi-homing device may also be vulnerable to spurious attacks via the different data interfaces available at the multi-homing device. For example, in a typical multi-homing device, an application installed therein can receive data from any network data interface as long as the data interface is open for data transfer and the data protocol addresses, e.g., IP address, port number, etc., match with that of the application.\nApart from security considerations, data network providers are also concerned about the billing and usage of various services and technologies available to the mobile telephones on the data networks. For example, there is a certain cost associated with each new service and technology that a data network provider provides and carriers are typically interested in hassle free discrete billing of various services used by the mobile phone user. If an application within a mobile telephone is restricted to use certain data interfaces available at the mobile telephone for data transfers, it can be easier for the carrier network to track the billing and cost associated with the usage of the different technologies and services distinctly, based on per-data interface usage.\nAdditionally, in a multi-homing networking device, the port space for networking applications is usually shared between all of the data interfaces available to the device. If an application is using a particular port number for data transfer on a particular data interface, no other application can use the same port number—even if the other application is using a completely different data interface. This can be an unnecessary restriction for a device that may need to run different services on different data interfaces but with the same port number. For example, the network device can include two different web servers that use the same port number, e.g., port eighty (80), but on different data interfaces. Most network devices do not allow this flexibility. Some implementations allow binding to a port for all of the data interfaces or for a specific data interface, i.e., one interface or all interfaces."}
{"text": "The brain consists of a plurality of neurons that communicate with each other via chemical messengers. Each neuron generates neurochemicals, referred to as neurotransmitters; and the neurotransmitters act at sites on the cellular membrane of neurons, the sites being referred to as receptors. One group of neurotransmitters, referred to as the monoamine neurotransmitters, includes serotonin, dopamine and noradrenaline.\nMonoamine neurotransmitters are released into the synaptic cleft in order to stimulate postsynaptic receptor activity. The removal (or inactivation) of monoamine neurotransmitters occurs mainly by a reuptake mechanism into presynaptic terminals. By inhibiting the re-uptake an enhancement of the physiological activity of monoamine neurotransmitters occur.\nNoradrenalin and serotonin re-uptake inhibitors are currently used as pharmaceuticals in anti-depressant therapy (Desipramine, Nortriptyline, and Protriptyline are inhibitors of noradrenaline-reuptake and Imipramine and Amitriptyline are mixed serotonine-reuptake and noradrenaline-reuptake inhibitors).\nThe pathophysiology of major affective illness is poorly understood, and several neurotransmitters have been implicated in the pathophysiology of major depression. However, several lines of preclinical and clinical evidence indicate that an enhancement of serotonin-mediated neurotransmission might underlie the therapeutic effect of the most recent and currently used drugs in anti-depressant therapy: Fluoxetine, Citalopram and Paroxetine.\nParadoxical serotonin re-uptake inhibitors inhibit the serotonin transporter within minutes whereas their full anti-depressant effect is seen only after three to four weeks of treatment, indicating that re-uptake inhibition per se is not responsible for the antidepressant response, but rather that further adaptive changes underlie and/or contribute to their therapeutic effect. The delayed onset of anti-depressant effect is considered to be a serious drawback to currently used monoamine re-uptake inhibitors.\nThe compounds provided with the present invention are potent serotonin (5-hydroxy-tryptamine, 5-HT) re-uptake inhibitors. The compounds of the invention also have noradrenaline and dopamine re-uptake inhibiting activity, the serotonin re-uptake and the noradrenaline re-uptake inhibiting activity of the compounds of the invention being stronger than the dopamine re-uptake inhibiting activity of the compounds, see table below.\nA strong dopamine re-uptake inhibiting activity is currently considered with the risk of undesirable central stimulating effects. On the other hand, an activating effect on the mesolimbic dopamine system is currently believed to underlay the commen mechanism of current antidepressant treatment by a mechanism which enhances the endogenous reward system. Compounds with a strong serotonin re-uptake inhibiting activity combined with a well balanced moderate dopamine re-uptake inhibiting activity may therefore provide agents with a rapid onset of anti-depressant effect.\nThe serotonergic neural system of the brain have been shown to influence a variety of physiologic functions, and the compounds of the present invention are predicted to have the ability to treat in mammals, including humans, a variety of disorders associated with this neural system, such as eating disorders, depression, obsessive compulsive disorders, panic disorders, alcoholism, pain, memory deficits and anxiety. Therefore, the present invention also provides methods of treating several disorders linked to decreased neurotransmission of serotonin in mammals. Included among these disorders are depression and related disorders such as pseudodementia or Ganser's syndrome, migraine pain, bulimia, obesity, pre-menstrual syndrome or late luteal phase syndrome, alcoholism, tobacco abuse, panic disorder, anxiety, post-traumatic syndrome, memory loss, dementia of ageing, social phobia, attention deficit hyperactivity disorder, chronic fatigue syndrome, premature ejaculation, erectile difficulty, anorexia nervosa, disorders of sleep, autism, mutism or trichotillomania.\nIn addition, compounds with dopamine re-uptake inhibiting activity are also considered useful for the treatment of Parkinsonism, depression, obesity, narcolepsy, drug addiction or misuse, attention-deficit hyperactivity disorders and senile dementia. Dopamine re-uptake inhibitors enhances indirectly via the dopamine neurones the release of acetylcholin and are therefore useful for the treatment of memory deficits, e.g. in Alzheimers disease and presenile dementia, and chronic fatigue syndrome. Noradrenaline re-uptake inhibitors are considered useful for enhancing attention, alertness, arousal, vigilance and for treating depression."}
{"text": "The present invention relates to a gas laser oscillator equipped with a constant laser gas pressure system which controls the laser gas pressure to be constant by exhausting and supplying certain amount of laser gas.\nFIG. 3 shows outline structure of a conventional gas laser oscillator including a constant laser gas pressure system. In the following, a conventional method for controlling the laser gas pressure constant in a gas laser oscillator is described with reference to FIG. 3.\nA conventional gas laser oscillator with constant laser gas pressure system comprises a power supply 1 for laser gas excitation by a discharge, a laser gas chamber 2, a vacuum pump 3, a laser gas bottle 4, an electromagnetic valve 5 for laser gas supply, a laser gas pressure sensor 6 for detecting the laser gas pressure within the laser gas chamber 2 and a laser gas pressure controller 7 which opens/closes (i.e performs open-and-close actions) the electromagnetic valve 5 for laser gas supply in order to maintain the pressure constant by comparing the signals detected by the laser gas pressure sensor 6 with a reference value.\nThe above-description illustrates structure of the constant laser gas pressure system of a conventional gas laser oscillator; in the following, operation of the system is described.\nAs a general rule among the gas laser oscillators, it is essential to keep laser gas pressure within the laser gas chamber 2 connected with the excitation power supply 1 constant, in order to maintain a stable laser output. Also, in order to prevent decrease of the decrease of output power due to deterioration of the laser gas, new laser gas is supplied continuously, or intermittently, from the laser gas bottle 4 or the like means connected with the gas laser oscillator.\nThe vacuum pump 3 connected to the laser gas chamber 2 exhausts deteriorated laser gas out of the laser gas chamber 2 as the routine work. When laser gas is exhausted the pressure of laser gas within the laser gas chamber 2 drops, and then the electromagnetic valve 5 for laser gas supply is opened to supply new laser gas into the laser gas chamber 2 from the laser gas bottle 4 connected to the gas laser oscillator. Quantity of the exhaust and that of the supply need to be well balanced; otherwise the pressure within the laser gas chamber 2 can not be kept constant. In order to keep the pressure constant, the laser gas pressure sensor 6 connected to the laser gas chamber 2 detects laser gas pressure within the laser gas chamber 2 and transmits signals of detected results to the laser gas pressure controller 7. The laser gas pressure controller 7 compares the signals of detected results with a reference value, or target value for the laser gas pressure, and opens/closes the electromagnetic valve 5 so that the laser gas pressure is maintained constant. The cycle period of open/close actions is kept stable when a gas laser oscillator is working normally.\nThe above-configured conventional gas laser oscillator has following drawbacks.\nAfter a long operation period, a gas laser oscillator sometimes exhibits a sudden drop in the laser output, or stops, due to leakage caused by a deteriorated vacuum seal or the like of the laser gas chamber 2 (hereinafter referred to as xe2x80x9cvacuum leakxe2x80x9d). In most of the cases, the vacuum leak does not take place abruptly, but it starts from a micro leak which hardly affects the laser output and it gradually develops into a leak that matters. Therefore, if the micro leak is detected at an early stage, a certain preventive maintenance measures can be taken to avoid such drop of laser output. Measurement of the vacuum leak has been conducted during the time when a gas laser oscillator is out of duty, using a measuring equipment. Lately, an increasing number of gas laser oscillators are put into operation continuously around the clock. On the other hand, there has been no effective means for detecting the micro vacuum leak applicable to gas laser oscillators at work. Once a sudden drop of laser output actually happens, scheduled maintenance is hardly helpful; in addition, the gas laser oscillator that is out of order is put aside for a considerable length of time before it is repaired. Thus, the inconvenience spreads over a production floor related to the processing with gas laser.\nThe present invention addresses the above-described problems.\nA gas laser oscillator of the present invention comprises:\na power supply for exciting laser gas by a discharge;\na laser gas chamber for containing the laser gas;\nmeans for exhausting laser gas contained in the laser gas chamber;\na gas pressure sensor for detecting gas pressure in the laser gas chamber;\nan electromagnetic valve for supplying laser gas into the laser gas chamber from a gas bottle so that pressure of the laser gas as detected by the gas pressure sensor is maintained constant;\ndetection means for detecting deviation in a cycle period of open/close actions of the electromagnetic valve from a reference value; and\nan alarm display controller for issuing an alarm based on a detection results delivered from the detection means.\nA gas laser oscillator of the present invention is also provided with a stop controller, which suspends the excitation of laser gas discharge in a case when the cycle period of open/close actions of the electromagnetic valve becomes longer than a reference value."}
{"text": "Field of the Invention\nThe present invention relates generally to defense apparatuses for domesticated animals. More particularly, the present invention relates to an improved device for a domesticated animal to wear to prevent prolonged animal fights.\nBackground\nIt is common for dogs to become aggressive with one another and engage in fights. Commonly, during a fight a dog will attempt to grab hold of its opponent's neck with its teeth in order to kill the other dog or force the other dog into submission. Dog fights pose dangers to both the animals involved and, unfortunately, to dog owners as they can be bitten in the attempt to stop the dogs from fighting. Dog fights can also cost hundreds or thousands of dollars in veterinarian bills for the dogs and medical bills for injured owners.\nSeveral methods are commonly attempted to get dogs to stop fighting. If one has access to a hose, spraying the dogs with water will often cause the dogs to stop fighting with one another. Likewise there are various chemicals which can be sprayed on the dogs which are irritating to the dogs and thereby cause the dogs to stop fighting. However, people often do not have access to a hose while walking their dogs or otherwise out in public with their dogs. Additionally, dog owners usually do not want to spray their dog with an irritant in case they encounter an aggressive dog, and getting close enough to spray the irritants on fighting dogs can put the dog owners at risk.\nThus, there is a need for an apparatus and method to assist pet owners in defending their animals from other aggressive animals. It may be advantageous for the apparatus or system to be simple to deploy and be available for use at all times when other animals are present."}
{"text": "The present invention relates to a manually guided implement, having a drive motor, and a grip for guiding the implement. A gas throttle is pivotably mounted on the grip and is connected by a transmitting means to a carburetor of the drive motor. An arresting element is held on the gas throttle for the releaseable arresting of the throttle on the grip, whereby the arresting element is adjustable relative to the gas throttle.\nU.S. Pat. No. 4,302,880 discloses a gas throttle having an arresting element for a motor-driven implement, namely a power chainsaw. The arresting element is rotatably mounted in the grip of the implement, and by means of finger-like dogs formed on one end thereof can be brought into engagement with the gas throttle, which is rotatably mounted on the grip. With the aid of the arresting element the gas throttle can be arrested, for example, in a starting throttle position. The arresting element is operatively connected with the gas throttle in such a way that operation of the gas throttle is precluded unless the arresting element is first actuated. Unfortunately, the structural configuration of this known arrangement is very complicated and unreliable during operation.\nIt is also known to embody a gas throttle of a manually guided motor-driven implement, such as a brush cutter, in such a way that it is monolithic with an arresting element, whereby the arresting element can be brought into arresting engagement with the grip of the implement by elastic deformation by the operator. In the arresting engagement, a starting throttle position of the gas throttle is provided. If the arresting position is unintentionally maintained over a long period of time, for example during the winter, the elastic restoration of the arresting element on the gas throttle is no longer ensured, for example due to material fatigue; the idling position of the gas throttle can then no longer be achieved.\nIt is therefore an object of the present invention to improve a manually guided implement of the aforementioned general type in such a way that the starting throttle position of the gas throttle can be easily set and can again be released in a functionally reliable manner."}
{"text": "1. Field of the Invention\nThe present invention relates generally to integrated circuit devices and, more specifically, to a flip-chip having an on-chip decoupling capacitor.\n2. Description of Related Art\nIt is generally desirable to have a semiconductor package which is more efficient and has high decoupling capacitance/low inductance. It is known that the effective inductance can be lowered by connecting decoupling capacitors nearer to a circuit (i.e, the effective inductance is lower by reducing the lead length). Inductance is a function of path length, therefore the longer the current path, the higher the inductance. High inductance, which yields higher supply noise in semiconductor packages, reduces the performance of integrated circuits (ICs). Also, inductance between an IC and power supply can induce spurious voltage spikes in the power supply system, which can in turn cause timing problems in signal switching.\nDecoupling capacitors are housed on semiconductor packages in order to lower the inductance through the package by reducing the lead length. Decoupling capacitors placed close to power consuming circuits are able to smooth out voltage variation with a stored charge on the decoupling capacitor. The stored charge either dissipates or is used as a local power supply to device inputs during signal switching stages, allowing the decoupling capacitor to negate the effects of voltage noise induced into the system by parasitic inductance. Off-chip decoupling capacitors, however, are not sufficient for very high speed microprocessor applications. Since the decoupling capacitors are located at a relatively long distance from the switching circuits, the time delay caused by the long inductance path makes the off-chip capacitors unusable with gigahertz switching circuits.\nIn order to sustain high frequency circuit operation, an ample amount of capacitive decoupling must be provided close to the switching circuits. Although it is possible to integrate chip capacitors within the chip's circuit elements, the capacitors compete for valuable die area that could be used for building additional circuits. Due to the limited area in which to build these capacitors, the overall capacitive decoupling that they provide is also limited."}
{"text": "Biosensors are used for the detection of molecules in biological samples, for instance proteins or DNA for diagnostic applications. It is also desired to detect drugs, therapeutic or abuse, in blood, urine or saliva. Such tests are developed to be used in many different settings and surroundings, e.g. at the point of care for medical applications, or at any desired place for drugs of abuse, e.g. at the roadside. In all cases a robust, reliable and sensitive device is required, which must also be low cost since it needs to be disposed after the measurement.\nCarrying out such a biochemical assay requires a certain degree of fluid handling, at least the sample fluid must be introduced in the sensing device in order to allow binding of the target molecules to the sensor surface. Depending on the kind of assay more or less complicated microfluidic systems are designed. Since the sample is contaminating it must not get in contact with the instrument and must be stored safely inside the cartridge during and after the measurement.\nRecently fully integrated micro fluidic on chip biochemical systems or lab on a chip systems have been developed. An issue in these micro fluidic systems is the manipulation of the fluids from and to the different reaction chambers, for which micro actuators such as pumps and valves are needed. Pumping and valving can be done in numerous ways. Depending on the application, for example type of assay, performance requirements, and cost requirements, the actuation of the fluid for dissolution of reagents, incubation, binding and washing, as an example, is implemented in different ways. There is a trade-off between the degree of control and simplicity, wherein simplicity can be identified with low costs. Either the fluid is actuated directly by mechanical metering with pistons or the fluid is not actuated but driven by capillary forces, so-called passive driving.\nThe latter is a cost-effective solution but does not allow for flow reversal, the flow rate is limited and not constant with distance and most importantly depending on the viscosity and surface tension of the fluid. A change in the required flow characteristics needs to be implemented in ‘hardware’ on the disposable which makes the system less flexible. Mechanical driving on the other hand is very flexible but requires physical contact which creates problems with operational lifetime of the instrument and contamination i.e. cleaning problems.\nTherefore there may be a need for an alternative low-cost fluid actuation technology in microfluidic systems in particular for medical disposables, like biosensors."}
{"text": "Because of the risk of communicable diseases, a number of syringes and adapters have been developed that are intended to prevent accidental needle sticks and/or inadvertent reuse of a syringe. The syringe safety systems come in a variety of forms. One form of a syringe safety system uses a needle guard comprising a body and an extendable shield. The shield can be manually or passively transitioned from a first, retracted position to a second, extended position following administration of the medicine within the syringe. In the extended position, the shield covers the needle, thereby protecting the user from accidental needle sticks.\nIn use, the syringe system is typically held by a user with two or more fingers engaging a finger grip of the syringe system and the user's thumb disposed on a thumb pad or top surface of the syringe's plunger. In many present systems, the finger grip area is quite small resulting in a device that may be difficult to manipulate and control, especially for end users with joint pain or limited dexterity.\nAccordingly, an improved finger grip area for use with safety syringe systems would be useful."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the alignment of a substrate in a lithographic projection apparatus. More specifically, it relates to alignment of the substrate after some process layers have been deposited above an alignment mark.\n2. Description of the Related Art\nThe term, “patterning means”, “patterning structure” or “mask” as here employed should be broadly interpreted as referring to means that can be used to endow an incoming radiation beam with a patterned cross-section, corresponding to a pattern that is to be created in a target portion of the substrate; the term “light valve” can also be used in this context. Generally, the said pattern will correspond to a particular functional layer in a device being created in the target portion, such as an integrated circuit or other device (see below). Examples of such patterning means include:\nA mask. The concept of a mask is well known in lithography, and it includes mask types such as binary, alternating phase-shift, and attenuated phase-shift, as well as various hybrid mask types. Placement of such a mask in the radiation beam causes selective transmission (in the case of a transmissive mask) or reflection (in the case of a reflective mask) of the radiation impinging on the mask, according to the pattern on the mask. In the case of a mask, the support structure will generally be a mask table, which ensures that the mask can be held at a desired position in the incoming radiation beam, and that it can be moved relative to the beam if so desired.\nA programmable mirror array. An example of such a device is a matrix-addressable surface having a viscoelastic control layer and a reflective surface. The basic principle behind such an apparatus is that (for example) addressed areas of the reflective surface reflect incident light as diffracted light, whereas unaddressed areas reflect incident light as undiffracted light. Using an appropriate filter, the said undiffracted light can be filtered out of the reflected beam, leaving only the diffracted light behind; in this manner, the beam becomes patterned according to the addressing pattern of the matrix-adressable surface. The required matrix addressing can be performed using suitable electronic means. More information on such mirror arrays can be gleaned, for example, from U.S. Pat. Nos. 5,296,891 and 5,523,193, which are incorporated herein by reference. In the case of a programmable mirror array, the said support structure may be embodied as a frame or table, for example, which may be fixed or movable as required.\nA programmable LCD array. An example of such a construction is given in U.S. Pat. No. 5,229,872, which is incorporated herein by reference. As above, the support structure in this case may be embodied as a frame or table, for example, which may be fixed or movable as required.\nFor purposes of simplicity, the rest of this text may, at certain locations, specifically direct itself to examples involving a mask and mask table; however, the general principles discussed in such instances should be seen in the broader context of the patterning means as hereabove set forth.\nLithographic projection apparatus can be used, for example, in the manufacture of integrated circuits (ICs). In such a case, the patterning means may generate a circuit pattern corresponding to an individual layer of the IC, and this pattern can be imaged onto a target portion (e.g. comprising one or more dies) on a substrate (silicon wafer) that has been coated with a layer of radiation-sensitive material (resist). In general, a single wafer will contain a whole network of adjacent target portions that are successively irradiated via the projection system, one at a time. In current apparatus, employing patterning by a mask on a mask table, a distinction can be made between two different types of machine. In one type of lithographic projection apparatus, each target portion is irradiated by exposing the entire mask pattern onto the target portion at once; such an apparatus is commonly referred to as a wafer stepper. In an alternative apparatus—commonly referred to as a step-and-scan apparatus—each target portion is irradiated by progressively scanning the mask pattern under the projection beam in a given reference direction (the “scanning” direction) while synchronously scanning the substrate table parallel or anti-parallel to this direction; since, in general, the projection system will have a magnification factor M (generally<1), the speed V at which the substrate table is scanned will be a factor M times that at which the mask table is scanned. More information with regard to lithographic devices as here described can be gleaned, for example, from U.S. Pat. No. 6,046,792, incorporated herein by reference.\nIn a manufacturing process using a lithographic projection apparatus, a pattern (e.g. in a mask) is imaged onto a substrate that is at least partially covered by a layer of radiation-sensitive material (resist). Prior to this imaging step, the substrate may undergo various procedures, such as priming, resist coating and a soft bake. After exposure, the substrate may be subjected to other procedures, such as a post-exposure bake (PEB), development, a hard bake and measurement/inspection of the imaged features. This array of procedures is used as a basis to pattern an individual layer of a device, e.g. an IC. Such a patterned layer may then undergo various processes such as etching, ion-implantation (doping), metallization, oxidation, chemo-mechanical polishing, etc., all intended to finish off an individual layer. If several layers are required, then the whole procedure, or a variant thereof, will have to be repeated for each new layer. Eventually, an array of devices will be present on the substrate (wafer). These devices are then separated from one another by a technique such as dicing or sawing, whence the individual devices can be mounted on a carrier, connected to pins, etc. Further information regarding such processes can be obtained, for example, from the book “Microchip Fabrication: A Practical Guide to Semiconductor Processing”, Third Edition, by Peter van Zant, McGraw Hill Publishing Co., 1997, ISBN 0-07-067250-4, incorporated herein by reference.\nFor the sake of simplicity, the projection system may hereinafter be referred to as the “lens”; however, this term should be broadly interpreted as encompassing various types of projection system, including refractive optics, reflective optics, and catadioptric systems, for example. The radiation system may also include components operating according to any of these design types for directing, shaping or controlling the projection beam of radiation, and such components may also be referred to below, collectively or singularly, as a “lens”. Further, the lithographic apparatus may be of a type having two or more substrate tables (and/or two or more mask tables). In such “multiple stage” devices the additional tables may be used in parallel, or preparatory steps may be carried out on one or more tables while one or more other tables are being used for exposures. Twin stage lithographic apparatus are described, for example, in U.S. Pat. No. 5,969,441 and WO 98/40791, incorporated herein by reference.\nA very important criterion in semiconductor manufactures is the accuracy with which the successive layers printed on the substrate are aligned with each other. Mis-alignments of the layers are referred to as overlay errors, for all the many layers required to make an integrated circuit the overlay errors must be kept within tight limits for the resulting device to function correctly. To correctly align the substrate to the mask and consequently minimize overlay errors, alignment marks, which generally take the form of diffraction gratings are etched in the bare silicon substrate. These alignment marks (referred to as “zero marks”) are aligned to corresponding marks provided on the mask using a variety of techniques, including through the lens (TTL) alignment systems and off-axis alignment systems. An example of the latter is described in EP-A-0 906 590 (P-0070). However, once a few process layers have been deposited or grown on the substrate, the zero marks etched in the bare substrate often become obscured and are no longer visible to the radiation used in the alignment process. Even if not completely obscured, the growth of layers on top of the alignment marks can be uneven, leading to a shift in the apparent position of the alignment mark. To enable alignment after the zero marks have been obscured, further alignment marks are printed during the deposition of suitable layers of the device. The subsequent marks, referred to as non-zero marks, are however subject to damage during subsequent process steps and will also accumulate overlay errors from previous process layers. When etching a blanket aluminum layer to define the interconnects of the integrated circuit, it is preferred to align to the original zero marks but to do this requires that the overlying aluminum layers, and possibly also dielectric layers, be removed. Such clearout steps are undesirable.\nA technique known as Impulsive Stimulated Thermal Scattering (ISTS) for measuring acoustic and thermal film properties, such as elastic constants and thermal diffusion rates, has been described in various publications such as J. A. Rogers et al., Appl. Phys. Lett 71 (2), 1997; A. R. Duggal et al. J. Appl. Phys. 72 (7), 1992; R. Logan et al., Mat. Res. Soc. Symp. Proc. 440, pg 347, 1997; L Dhar et al., J. Appl. Phys. 77 (9), 1995; and J. A. Rogers et al. Physica B 219 & 220, 1996. In this method, two excitation pulses overlapping in time and space are incident on a sample at slightly different angles. The two pulses interfere and heat the sample in a pattern corresponding to the interference pattern between them. The local heating sets up vibrations in the crystal structure of the sample which act as a diffraction grating to a probe pulse incident on the sample shortly after the excitation pulses. The diffraction of the excitation pulse is measured to give an indication of the property being investigated in the sample."}
{"text": "For example, in the case where a concentration of moisture contained in the exhaust gas is measured by a non-dispersive infrared absorption (NDIR) method, the exhaust gas is cooled and the moisture is condensed, and this results in an erroneous measurement value different from that of the original one. In order to prevent such an erroneous measurement, a sampling line up to reach the moisture concentration meter is heated in some cases (JP2002-340778A).\nIn this configuration, it is configured that the temperature of the exhaust gas in the sampling line is regulated to be equal to or higher than 100° C. to thereby prevent the moisture contained in the exhaust gas from being condensed. Further, even in the case where only the concentration other than the moisture contained in the exhaust gas is measured, when the concentration of the moisture is largely changed, the measurement accuracy of the other components is affected. Therefore, it is configured so that the temperature of the exhaust gas in the sampling line is also regulated to be equal to or higher than 100° C.\nHowever, since the vehicle mounted type exhaust gas analyzing apparatus uses, for example, a battery of the vehicle as a power supply source thereof and the power capacity is limited. Therefore, in the case where a large power consumption occurs in such a temperature regulation mechanism, an operational time of the exhaust gas analyzing apparatus is extremely limited and it is difficult to perform a sufficient exhaust gas analysis."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication terminal apparatus to be used for a remote meeting (a meeting held between persons who are physically apart from each other, a certain communication system thus being needed to enable these persons to communicate with each other), which apparatus has such functions as communicating through so-called tele-writing and controlling documents.\n2. Discussion of the Background\nA tele-writing apparatus is defined in the CCITT (International Telegraph and Telephone Consultative Committee) Recommendation T.150, as being a remote communication terminal which communicates with another remote communication terminal via a communication line, which terminal has a display device and a hand writing inputting device, which terminal transfers hand writing input information, that is, information which has been input by hand writing, to the other terminal while a conversation is taking place, and which terminal helps the persons taking part in the conversation to understand the transferred hand writing input information by displaying it.\nThe Shimazu corporation has developed such a tele-writing apparatus named \"pen-telephone\" (type TW101, hand writing communication terminal equipment with a facsimile function), which apparatus displays, on respective liquid crystal displays of both the present apparatus and another apparatus, drawing information comprising characters and diagrams, the liquid crystal displays having transparent tablets respectively formed therein, and inputs information by hand writing on the transparent tablets so that the corresponding input image information is displayed on the liquid crystal displays of both the apparatuses in real time.\nFurther, in a system for a remote meeting utilizing such tele-writing apparatuses, it is necessary to communicate, that is, transmit and receive, and treat the following two types of data: a bulk type data and a real time type data. The bulk type data is typically utilized as a facsimile data based on raster images, and the real time type data is typically utilized as a tele-writing data based on position information.\nEither only the tele-writing apparatus itself may be used, or the tele-writing apparatus may be used together with a display of a raster image, in a remote meeting. In the remote meeting, a common document is displayed on both apparatuses, the contents of the displayed document being then controlled (modified) through the real time type data. The document controlling (modifying) may include specifying the document, modifying it by hand writing, enlarging/reducing the size thereof, scrolling the contents, and rotating a direction thereof.\nIn the \"pen-telephone\", it is possible for the hand writing inputting to be performed only with respect to the presently displayed contents. Thus, in a case where a memorandum to be referred to in the present conversation includes a plurality of pages, it is necessary that the presently displayed contents is to be deleted after the contents are printed out once, and further, it is then necessary that the other hand written information be written.\nThis action interrupts the smooth proceeding of a remote meeting, and further the previous drawing information, which has been deleted after printing of the displayed contents, cannot be displayed again. Further, the \"pen-telephone\" does not have functions for storing a memorandum including a plurality of pages.\nFurther, the tele-writing apparatus according to the CCITT Recommendation T.150 only transmits/receives hand written information. It does not define other functions for transmitting/receiving a document and other functions for displaying/controlling contents of the document. Such a tele-writing apparatus cannot provide a smooth remote meeting.\nFurther, according to the Recommendation T.150, it is normal for the tele-writing apparatus to be used alone. Thus, it is possible to use the apparatus for making a memorandum of a meeting. Further, in a case where the operation mode of the tele-writing apparatus is changed between a first simple tele-writing mode and a second mode in which a document is transmitted/received, it is necessary to store the presently displayed file of information, to finish the current operation according to the first mode, to start another operation according to the second mode, and then to load the file. Such troublesome operations may interfere with the smooth proceeding of the meeting."}
{"text": "Depression has been a major social problem. A survey in the United States of America estimates that the lifetime prevalence rate of depressive disorders is 26% for women and 12% for men (Non-patent Document 1). An extensive survey was also carried out in Japan, and it was reported that the lifetime prevalence rate was 14%. It was also reported that 84% of the patients with depression have complication, 61% of the patients have other mental diseases, 30% of the patients have personality disorders, and 58% of the patients have complication of physical diseases (Non-patent Document 2).\nCommon symptoms of depression include a wide variety of symptoms, such as depressive mood, lethargy, depressive attitudes, suicidal ideation, impatience, insomnia, anorexia, decreased sexual impulse, and physical ailments. Pathologically-severe depression is called major depression. With regard to major depression, there is reported experimental data showing that in many cases, a rapid recovery is observed during the first 40 weeks, but the recovery plateaus thereafter. It is reported that results of a 20-year follow-up survey on patients having major depression show that 15 to 19% of the patients had chronic residual symptoms such as working difficulty (Non-patent Document 3).\nWith regard to causes of depression, various theories are suggested. A hypothesis that was suggested in the early days and is still popular is a theory stating that a decrease in the function of monoamines in the brain is responsible for depression (this hypothesis is called “intracerebral monoamine hypothesis”). This theory is suggested on the basis of the fact that a tricyclic antidepressive agent having an action to enhance the function of monoamines in the brain by blocking the reuptake of monoamines in cells is effective for remission of depression.\nCurrently, especially the function of serotonin, which is an intracerebral monoamine, attracts great attention, and there is a popular hypothesis suggesting that the antidepressive action can be produced by sensitizing serotonin-1 receptors present in the serotonin nerve cell body to increase the amount of free serotonin from presynaptic cells (Non-patent Document 4). This idea is basically the same as that of the classical intracerebral monoamine hypothesis. In any cases, it is considered that depression is caused by a decrease in the intracerebral serotonin function.\nWith regard to medications for depression, first, a tricyclic antidepressive agent, imipramine, and monoamine oxidase inhibitors were developed about 40 years ago. On the basis of the idea that the action to increase serotonin and noradrenaline in the synaptic cleft, which is one of the features the above drugs, is the core of the antidepressant effect, many antidepressive agents have been developed following imipramine and monoamine oxidase inhibitors. In recent years, use of selective serotonin reuptake inhibitors (SSRI) and serotonin/norepinephrine reuptake inhibitors (SNRI), which more selectively inhibit the reuptake of the monoamines by serotonin transporters, have been used in clinical practice. There are as many as 17 antidepressive agents that are approved in Japan at this time, including tricyclic/tetracyclic antidepressive agents, SSRI, and SNRI (Non-patent Document 5).\nThese antidepressive agents are effective against various diseases besides depression. According to package inserts of the pharmaceutical agents presented by the Pharmaceuticals and Medical Devices Agency, for example, fluvoxamine maleate, which is an SSRI, is effective against obsessive-compulsive and social anxiety disorders; paroxetine hydrochloride, which is an SSRI, is effective against panic and obsessive-compulsive disorders; and sertraline hydrochloride, which is an SSRI, is effective against a panic disorder. Further, amitriptyline hydrochloride, which is a tricyclic SNRI, is effective against nocturnal enuresis.\nHowever, it is reported that these antidepressive agents have side effects. For example, nausea, headache, hypersensitivity and the like are reported as side effects of SSRI, and tremor, tachycardia, erection/ejaculation disorders and the like are reported as side effects of SNRI. There are several reports that side effects are increased in the case of administration of multiple antidepressive agents, compared with administration of a single antidepressive agent (Non-patent Document 6).\nIt is also reported that when administration of an antidepressive agent is discontinued, withdrawal symptoms may occur. Withdrawal symptoms of tricyclic/tetracyclic antidepressive agents that are reported include: digestive symptoms and physical discomfort accompanied by anxiety and irritability, such as lethargy, emesis, and headache; sleep disorders such as insomnia and multi-dream; motility disorders such as akathisia and Parkinsonian symptoms; activation of behavior that leads to a transition to a manic state; and arrhythmia. These symptoms, except for arrhythmia, are also observed in withdrawal symptoms of SSRI (Non-patent Document 7). Disturbance to the sense of equilibrium, paresthesia, aggressive/impulsive behavior and the like are reported as withdrawal symptoms characteristic of SSRI (Non-patent Document 8).\nAs described above, it is reported that the antidepressive agents that are widely used at this time have side effects, and the development of safe food ingredients and components, which have ever been eaten and are applicable in place of the above antidepressive agents, are expected. Royal jelly-derived compositions (Patent Document 1) and compositions comprising a hop extract as a main component (Patent Document 2) are reported as food-derived compositions having antidepressive action.\n1-Methyl-1,2,3,4-tetrahydro-β-carboline-3-carboxylic acid (MTCA) is a compound contained in fruits, such as lemon, grapefruit, orange, and mandarin (Non-patent Document 9), beer, wine (Non-patent Document 10), soy sauce (Non-patent Document 11) and the like, and two types of diastereomers, (1S,3S) and (1R,3S), are known. It is known that MTCA is produced nonfermentatively in the presence of tryptophan and acetaldehyde; especially, it is known that (1S,3S)-MTCA (also called (1S,3S)-2,3,4,9-tetrahydro-1-methyl-1H-pyrido[3,4-b]-indole-3-carboxylic acid) is produced through fermentation by S. cerevisiae (Non-patent Document 12). Further, it is known that (1S,3S)-MTCA has antioxidation action, prevents self-polymerization reaction of γ-crystallin by the Fenton reaction using FeCl3, and prevents photopolymerization reaction of γ-crystallin (Non-patent Document 13). Further, there is a description that pinoline (6-methoxy-1,2,3,4-tetrahydro-β-carboline), which is a carboline compound, inhibits the activity of monoamine oxidase A, and also inhibits the uptake of serotonin in the brain (Non-patent Document 14). Further, it is reported that pinoline reduces immobility time in forced-swimming tests using rats, and that this is the same action as that of a tricyclic antidepressive agent (Non-patent Document 15).\nHowever, inhibitory action of serotonin uptake, antidepressive effect, and improvement of learning motivation by MTCA have completely been unknown."}
{"text": "In recent history, relatively inexpensive electronic devices have become available to ordinary property-owners that wish to monitor and control various aspects of their properties. A recent concept called the “Internet of Things” imagines home-related electronic devices that can be reached through the Internet, by which an environment can be controlled, e.g., lighting, temperature, digital video recorders, and many other “smart” devices. That kind of device ordinarily requires a connection to a network switch or hub, which connection can be wired or wireless.\nWireless collections to such smart devices are often desired, particularly in existing constructions, avoiding the laying of wires in existing walls and other structures. Technologies serving this purpose include wattage devices that communicate using the 2.4 GHz 802.11b/g “WiFi” protocol, and other more-recent and similar protocols such as Zigbee and Z-Wave. These protocols generally permit data rates of 100 k bytes per second or more, allowing for devices that transmit and forward audio and video data in substantial real-time. However with high data rates come a vulnerability to interference from other devices operating on the same radio bands, and accordingly devices using these short-range protocols are generally limited to service within a single residence or building within a distance of less than 100 meters.\nRecent technologies have been developed that permit operation to an intermediate range, communicating between points that are several miles or more away, for example using the LoRaWAN protocol. In this type of network, interference reduction is achieved by using frequencies in the UHF band and by including redundancies in communication, using for example multiple sampling, multiple frequency (spread-spectrum) techniques, and/or error-tolerant protocols. The use of the UHF band avoids interference from over-the-horizon sources, while at the same time avoiding some attenuation-of-signal from water-vapor, precipitation, buildings and other physical obstructions. These redundancies and protocols necessarily reduce the data throughput such that audio and video data cannot be streamed in good quality or in real-time.\nAn exemplary use of intermediate-range communication is in the recent deployment of wireless utility meters. Having a utility meter that can be read without a person traveling to and visually looking at it is a substantial cost savings for a utility. For such a use a meter communicates two items of information, which are an identifier for the meter and the meter reading itself; the utility takes a pair of such items and generates a bill for a utility subscriber. Because utility bills are ordinarily generated once per month, the amount of data from a single meter is usually on the order of a few tens of bytes in that period. Thus tens or even hundreds of thousands of meters can share a single intermediate-range channel, largely without interference from other devices.\nThe unsuitability of existing systems at intermediate ranges for large numbers of devices, for example in the Internet-Of-Things, is an unsolved problem. Turning now to FIG. 2, two proximal short-range networks are conceptually shown, such as those constructed from devices communicating through the 802.11 g protocol, for example using the Zigbee specification. In the example two property owners each operate a hub “H” servicing multiple devices “D”, and as shown in the figure the circumference of effective communication 10 of each hub is proximal so as to create a zone of interference 11. Communication with the devices in that zone can be affected by communications with a non-subscribing hub, so as to introduce possible packet transmission collisions and interference. However, because these networks are short-range, the zone of interference 11 is relatively small, and the frequency of collisions and interference is likewise negligible. Proximal operators of such networks can experience impacts, but because of the short-range and the short transmission time of packets, these impacts are minor and usually acceptable.\nIn contrast as shown in FIG. 3, three intermediate-range networks are located in similar proximity, each with a hub “H” and subscribing devices “D”. Because of the greater distance of communication provided, most devices “D” are located in zones of interference 11, and many devices may be located in regions 12 where several hubs are located within range. Thus where intermediate-range networks are to be used most communication between hubs and devices and also inter-device communications should be expected to take place under conditions of interference, especially when located within a city or other populated area. Utility meter reading and other existing installations using the LoRa or LoRaWAN protocols manage this congestion by being the only network in the radio locality on a particular frequency, and by infrequent packet transmission on the order of once per week or month such that collisions between devices aren't likely. Utility meters are configured and installed by the utility provider, who can limit the number of end-devices within a locality to ensure network reliability. Because the provider is the only consumer of the applicable bandwidth resources, it can effectively control interference and manage congestion.\nUnavailable on the consumer market today are hubs and end-devices that can operate at intermediate ranges. The short-range “WiFi” 802.11 b/g network is now so popular that in most urban areas there is a router within communication range no matter where a person finds himself; most households using the Internet now have one, purchased from a retailer. If the range of those routers and the devices they service were to be extended to an intermediate range, the result would be a collapse of system functionality from overwhelming interference. Thus, in order to bring intermediate-range devices to the ordinary consumer, further development is needed.\nThe existing intermediate-range techniques, however, aren't conducive for applications where interactivity is need. For a channel sharing thousands of meters, it isn't necessary to resolve collisions between devices in a matter of milliseconds, because data transmissions can be delayed without significant impacts. In another example, an irrigation controller will ordinarily keep a set of sprinklers on for minutes at a time, and a delay of multiple seconds or even minutes is tolerable. In contrast, a person activating a light switch, for example, will not accept activation of lights with perhaps more than a one-second delay. Where a person enters a code on a keypad to enter a building, he expects a controlled lock to deactivate in real-time. In general, the existing intermediate-range technologies are fault-susceptible and not reliable for such interactivity, particularly where multiple devices share a common communications frequency or channel.\nInteractivity issues for battery-powered devices can be even worse. For these devices, it is generally undesirable to keep a receiver continuously powered, and worse to repeatedly being awakened from a sleep mode to process and discriminate packets destined for other devices. The LoRaWAN Class A and B protocols address this by having end-devices turn off their receivers for long periods of time, waking up periodically to interact with a network gateway. Such a device may remain asleep for seconds, minutes or even hours, and thus cannot be made responsive to incoming queries. Furthermore, these protocols are susceptible to collisions from co-transmitting devices, which may require backing off interactions with a hub, and no time of reception can be guaranteed. Thus absent from the field of the invention is a system that can provide adequate and reliable service for groups of sensed and controlled remote devices at intermediate ranges.\nDisclosed herein are hub-based wireless networks employing end-devices at intermediate wireless at ranges of thousands of meters. End-devices may be such things as keypads, door latches, occupancy monitors, sprinkler controllers and other devices controlled or monitored in a campus or a collection of buildings. A device-initiated asynchronous mode may be used to communicate information related to spontaneous events to a hub. The hub receives packet retrial information from the end-devices, determines whether packet retrials are due to interference or some other cause such as collisions, and applies varying levels of redundancy based upon a recent history of success and failure of packet exchanges, which measurement may be a value of badness increased as errors are encountered and decreased as packet communication is successful. Detailed information on various example embodiments of the inventions are provided in the Detailed Description below, and the inventions are defined by the appended claims."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for processing a linework image, and more particularly to a method and apparatus for finding and reconnecting a break in a linework image.\n2. Description of the Related Art\nA block copy is prepared as an original of characters and lineworks in a prepress process for a color printing plate. The block copy is formed by arranging phototype-setting characters and drafted keylines on a layout sheet in the same dimensions and quality as the final products. Instructions for the subsequent processes are also given on the layout sheet; that is, the layout sheet acts as an instruction paper for the prepress process.\nTint laying is generally performed in the prepress process. Tint laying is processing in which a specified region within an image is uniformly filled with a desired color. Automatic tint laying is employed in some of the modern image processing systems. In such systems, tint laying is completed, for example, by obtaining binary image data of a block copy image with an image scanner and filling a certain closed region in the block copy image with a desired color with an image processor.\nA linework, which is supposed to define a closed region in a block copy image, is sometimes found incomplete or broken when filling of the closed region is executed. The incomplete or broken linework is generally attributable to an incomplete draft or layout sheet, or unsuccessful image reading with an image scanner.\nLine breaks are generally reconnected manually; an operator directly detects line breaks in a block copy image displayed on a CRT and fills them to correct the block copy image.\nA line break has generally a width of one to several pixels. Accordingly, only experienced operators can efficiently detect breaks on the CRT. When there are several breaks in a linework, it takes rather a long time to detect and reconnect all of them.\nJAPANESE PATENT LAYING-OPEN GAZETTE No. Sho-61-139892 proposes a method of reconnecting line breaks which utilizes a technique of changing binary image data into vector data. Vector data processors for executing this method are rather complicated and expensive.\nOn the other hand, those image processors for processing lineworks (hereinafter referred to as linework processors), which change the level of binary image data for each pixel, are relatively simple. Accordingly, addition of a specific function of vector data processing to linework processors drastically increases the price of the processors.\nFurthermore, the method including vector data processing does not perfectly detect line breaks and requires operators confirmation for the existence thereof. It means that the total time for processing is not shortened greatly."}
{"text": "1. Field of the Invention\nThis invention is in the field of devices for producing concrete floors.\n2. Description of the Prior Art\nConcrete slabs are used in animal pens to form slotted floors. The slabs are spaced apart to allow for the animal droppings to pass therethrough to a pit or ditch located beneath the floor. It is desirable for the concrete slabs to have a smooth top surface thereby allowing for the easy removal of the animal droppings from the slabs. In order to obtain the smooth top surface for each slab, I have provided a mold disclosed in my co-pending U.S. patent application, Ser. No. 457,068, now U.S. Pat. No. 3,915,422 wherein the slabs are poured upside-down. The mold is provided with a smooth bottom surface thereby forming a smooth top surface of the slabs which are simultaneously poured into an integral section.\nHeretofore, a floor section composed of a plurality of slabs has been formed by manually filling and handling the mold. This method limits the number of floor sections pourable in a single day and provides for inefficiency. I have therefore provided a system for pouring a plurality of molds. The system disclosed herein allows for a far greater number of floor sections to be produced in a single day.\nThe molds are provided with open tops and as a result, the mold must first be turned over in order to remove the mold from the poured floor section. The mold when filled with concrete is exceptionally heavy therefore requiring a hoist or other device for lifting and turning the mold over. Disclosed herein is an automatic system for turning the mold over as the mold progresses down a conveyor line. Likewise, it is necessary to turn the mold right-side-up once the mold has been removed from the floor section. Disclosed herein is a return conveyor which rotates the mold to a right-side-up position. A shaker table is provided allowing for settling of the concrete within the mold."}
{"text": "This invention relates to polynucleotides, (herein referred to as xe2x80x9cBASB023 polynucleotide(s)xe2x80x9d), polypeptides encoded by them (referred to herein as xe2x80x9cBASBO23xe2x80x9d or xe2x80x9cBASB023 polypeptide(s)xe2x80x9d). recombinant materials and methods for their production. In another aspect, the invention relates to methods for using such polypeptides and polynucleotides, including vaccines against bacterial infections. In a further aspect, the invention relates to diagnostic assays for detecting infection of certain pathogens.\nMoraxella catarrhalis (also named Branhamella catarrhalis) is a Gram negative bacteria frequently isolated from the human upper respiratory tract. It is responsible for several pathologies the main ones being otitis media in infants and children, and pneumonia in elderlies. It is also responsible of sinusitis, nosocomial infections and less frequently of invasive diseases.\nOtitis media is an important childhood disease both by the number of cases and its potential sequelae. More than 3.5 millions cases are recorded every year in the United States, and it is estimated that 80% of the children have experienced at least one episode of otitis before reaching the age of 3 (Klein, JO (1994) Clin.Inf.Dis 19:823). Left untreated, or becoming chronic, this disease may lead to hearing losses that could be temporary (in the case of fluid accumulation in the middle ear) or permanent (if the auditive nerve is damaged). In infants, such hearing losses may be responsible for a delayed speech learning.\nThree bacterial species are primarily isolated from the middle ear of children with otitis media: Streptococcus pneumoniae, non typeable Haemophilus influenza (NTHi) and M. catarrhalis. They are present in 60 to 90% of the cases. A review of recent studies shows that S. pneumoniae and NTHi represent both about 30%. and M. catarrhalis about 15% of the otitis media cases (Murphy. TF (1996) Microbiol.Rev. 60:267). Other bacteria could be isolated from the middle ear (H. influenza type B, S. pyogenes etc) but at a much lower frequency (2% of the cases or less).\nEpidemiological data indicate that, for the pathogens found in the middle ear, the colonization of the upper respiratory tract is an absolute prerequisite for the development of an otitis: other are however also required to lead to the disease (Dickinson, D P et al. (1988) J. Infect.Dis. 158:205, Faden, H L et al. (1991) Ann.Otorhinol.Laryngol. 100:612). These are important to trigger the migration of the bacteria into the middle ear via the Eustachian tubes, followed by the initiation of an inflammatory process. These factors are unknown todate. It has been postulated that a transient anomaly of the immune system following a viral infection, for example, could cause an inability to control the colonization of the respiratory tract (Faden, H L et al (1994) J. Infect.Dis. 169:1312). An alternative explanation is that the exposure to environmental factors allow a more important colonization of some children, who subsequently become susceptible to the development of otitis media because of the sustained presence of middle ear pathogens (Murphy, T F (1996) Microbiol.Rev. 60:267).\nThe immune response to M. catarrhalis is poorly characterized. The analysis of strains isolated sequentially from the nasopharynx of babies followed from 0 to 2 years of age, indicates that they get and eliminate frequently new strains. This indicates that an efficacious immune response against this bacteria is mounted by the colonized children (Faden, H L et al (19.94) J. Infect.Dis. 169:1312).\nIn most adults tested, bactericidal antibodies have been identified (Chapman, A J et al. (1985) J. Infect.Dis. 151:878). Strains of M. catarrhalis present variations in their capacity to resist serum bactericidal activity: in general, isolates from diseased individuals are more resistant than those who are simply colonized (Hol, C et al. (1993) Lancet 341:1281, Jordan, K L et al. (1990) Am.J.Med. 88 (suppl. 5A):28S). Serum resistance could therfore be considered as a virulence factor of the bacteria. An opsonizing activity has been observed in the sera of children recovering from otitis media.\nThe antigens targetted by these different immune responses in humans have not been identified, with the exception of OMIP B4, a 84 kDa protein which expression is regulated by iron, and that is recognized by the sera of patients with pneumonia (Sethi, S. et al. (1 995) Infect.Immun. 63:1516), and of UspA1 and UspA2 (Chen D. et al.(1999), Infect.Immun. 67:1310).\nA few other membrane proteins present on the surface of M. catarrhalis have been characterized using biochemical method, or for their potential implication in the induction of a protective immunity (for review, see Murphy, T F (1996) Microbiol.Rev. 60:267). In a mouse pneumonia model, the presence of antibodies raised against some of them (UspA, CopB) favors a faster clearance of the pulmonary infection. Another polypeptide (OMP CD) is highly conserved among M. catarrhalis strains, and presents homologies with a porin of Pseudomonas aeruginosa, which has been demonstrated efficacious against this bacterium in animal models.\nThe frequency of Moraxella catarrhalis infections has risen dramatically in the past few decades. This has been attributed to the emergence of multiply antibiotic resistant strains and an increasing population of people with weakened immune systems. It is no longer uncommon to isolate Moraxella catarrhalis strains that are resistant to some or all of the standard antibiotics. This phenomenon has created an unmet medical need and demand for new anti-microbial agents, vaccines, drug screening methods, and diagnostic tests for this organism.\nThe present invention relates to BASB023, in particular BASB023 polypeptides and BASB023 polynucleotides, recombinant materials and methods for their production. In another aspect, the invention relates to methods for using such polypeptides and polynucleotides, including prevention and treatment of microbial diseases, amongst others. In a further aspect, the invention relates to diagnostic assays for detecting diseases associated with microbial infections and conditions associated with such infections, such as assays for detecting expression or activity of BASB023 polynucleotides or polypeptides.\nVarious changes and modifications within the spirit and scope of the disclosed invention will become readily apparent to those skilled in the art from reading the following descriptions and from reading the other parts of the present disclosure.\nThe invention relates to BASB023 polypeptides and polynucleotides as described in greater detail below. In particular, the invention relates to polypeptides and polynucleotides of BASB023 of Moraxella catarrhalis, which is related by amino acid sequence homology to Legionella adelaidensis macrophage infectivity potentiator polypeptide. The invention relates especially to BASB023 having the nucleotide and amino acid sequences set out in SEQ ID NO:1 or 3 and SEQ ID NO:2 or 4 respectively. It is understood that sequences recited in the Sequence Listing below as xe2x80x9cDNAxe2x80x9d represent an exemplification of one embodiment of the invention, since those of ordinary skill will recognize that such sequences can be usefully employed in polynucleotides in general, including ribopolynucleotides.\nPolypeptides\nIn one aspect of the invention there are provided polypeptides of Moraxella catarrhalis referred to herein as xe2x80x9cBASB023xe2x80x9d and xe2x80x9cBASB023 polypeptidesxe2x80x9d as well as biologically, diagnostically, prophylactically, clinically or therapeutically useful variants thereof, and compositions comprising the same.\nThe present invention further provides for:\n(a) an isolated polypeptide which comprises an amino acid sequence which has at least 85% identity, preferably at least 90% identity, more preferably at least 95% identity, most preferably at least 97-99% or exact identity, to that of SEQ ID NO:2 or 4;\n(b) a polypeptide encoded by an isolated polynucleotide comprising a polynucleotide sequence which has at least 85% identity, preferably at least 90% identity, more preferably at least 95% identity, even more preferably at least 97-99% or exact identity to SEQ ID NO:1 or 3 over the entire length of SEQ ID NO:1 or 3 respectively; or\n(c) a polypeptide encoded by an isolated polynucleotide comprising a polynucleotide sequence encoding a polypeptide which has at least 85% identity, preferably at least 90% identity, more preferably at least 95% identity, even more preferably at least 97-99% or exact identity, to the amino acid sequence of SEQ ID NO:2 or 4.\nThe BASB023 polypeptides provided in SEQ ID NO:2 or 4 are the BASB023 polypeptides from Moraxella catarrhalis strains Mc293 I (ATCC 43617).\nThe invention also provides an immunogenic fragment of a BASB023 polypeptide, that is, a contiguous portion of the BASB023 polypeptide which has the same or substantially the same immunogenic activity as the polypeptide comprising the amino acid sequence of SEQ ID NO:2 or 4; That is to say, the fragment (if necessary when coupled to a carrier) is capable of raising an immune response which recognises the BASB023 polypeptide. Such an immunogenic fragment may include, for example, the BASB023 polypeptide lacking an N-terminal leader sequence, and/or a transmembrane domain and/or a C-terminal anchor domain. In a preferred aspect the immunogenic fragment of BASB023 according to the invention comprises substantially all of the extracellular domain of a polypeptide which has at least 85% identity, preferably at least 90% identity, more preferably at least 95% identity, most preferably at least 97-99% identity, to that of SEQ ID NO:2 or 4 over the entire length of SEQ ID NO:2 or 4.\nA fragment is a polypeptide having an amino acid sequence that is entirely the same as part but not all of any amino acid sequence of any polypeptide of the invention. As with BASB023 polypeptides, fragments may be xe2x80x9cfree-standing,xe2x80x9d or comprised within a larger polypeptide of which they form a part or region, most preferably as a single continuous region in a single larger polypeptide.\nPreferred fragments include, for example, truncation polypeptides having a portion of an amino acid sequence of SEQ ID NO:2 or 4 or of variants thereof, such as a continuous series of residues that includes an amino- and/or carboxyl-terminal amino acid sequence. Degradation forms of the polypeptides of the invention produced by or in a host cell, are also preferred. Further preferred are fragments characterized by structural or functional attributes such as fragments that comprise alpha-helix and alpha-helix forming regions, beta-sheet and beta-sheet-forming regions, turn and turn-forming regions, coil and coil-forming regions, hydrophilic regions, hydrophobic regions, alpha amphipathic regions, beta amphipathic regions, flexible regions, surface-forming regions, substrate binding region. and high antigenic index regions.\nFurther preferred fragments include an isolated polypeptide comprising an amino acid sequence having at least 15, 20, 30, 40, 50 or 100 contiguous amino acids from the amino acid sequence of SEQ ID NO:2 or 4, or an isolated polypeptide comprising an amino acid sequence having at least 15, 20, 30, 40, 50 or 100 contiguous amino acids truncated or deleted from the amino acid sequence of SEQ ID NO:2 or 4.\nFragments of the polypeptides of the invention may be employed for producing the corresponding full-length polypeptide by peptide synthesis; therefore, these fragments may be employed as intemiediates for producing the full-length polypeptides of the invention.\nParticularly preferred are variants in which several, 5-10, 1-5, 1-3, 1-2 or 1 amino acids are substituted, deleted, or added in any combination.\nThe polypeptides, or immunogenic fragments, of the invention may be in the form of the xe2x80x9cmaturexe2x80x9d protein or may be a part of a larger protein such as a precursor or a fusion protein. It is often advantageous to include an additional amino acid sequence which contains secretory or leader sequences, pro-sequences, sequences which aid in purification such as multiple histidine residues, or an additional sequence for stability during recombinant production. Furthermore, addition of exogenous polypeptide or lipid tail or polynucleotide sequences to increase the immunogenic potential of the final molecule is also considered.\nIn one aspect, the invention relates to genetically engineered soluble fusion proteins comprising a polypeptide of the present invention, or a fragment thereof, and various portions of the constant regions of heavy or light chains of immunoglobulins of various subclasses (IgG, IgM, IgA, IgE). Preferred as an immunoglobulin is the constant part of the heavy chain of human IgG, particularly IgG1, where fusion takes place at the hinge region. In a particular embodiment, the Fc part can be removed simply by incorporation of a cleavage sequence which can be cleaved with blood clotting factor Xa.\nFurthermore, this invention relates to processes for the preparation of these fusion proteins by genetic engineering, and to the use thereof for drug screening, diagnosis and therapy. A further aspect of the invention also relates to polynucleotides encoding such fusion proteins. Examples of fusion protein technology can be found in International\nThe proteins may be chemically conjugated, or expressed as recombinant fusion proteins allowing increased levels to be produced in an expression system as compared to non-fused protein. The fusion partner may assist in providing T helper epitopes (immunological fusion partner), preferably T helper epitopes recognised by humans, or assist in expressing the protein (expression enhancer) at higher yields than the native recombinant protein. Preferably the fusion partner will be both an immunological fusion partner and expression enhancing partner.\nFusion partners include protein D from Haemophilus influenzae and the non-structural protein from influenzae virus, NS1 (hemagglutinin). Another fusion partner is the protein known as LytA. Preferably the C terminal portion of the molecule is used. Lyta is derived from Streptococcus pneumoniae which synthesize an N-acetyl-L-alanine amidase, amidase LytA, (coded by the lytA gene {Gene, 43 (1986) page 265-272}) an autolysin that specifically degrades certain bonds in the peptidoglycan backbone. The C-terminal domain of the LytA protein is responsible for the affinity to the choline or to some choline analogues such as DEAE. This property has been exploited for the development of E.coli C-LytA expressing plasmids useful for expression of fusion proteins. Purification of hybrid proteins containing the C-LytA fragment at its amino terminus has been described {Biotechnology: 10, (1992) page 795-798}. It is possible to use the repeat portion of the LytA molecule found in the C terminal end starting at residue 178, for example residues 188-305.\nThe present invention also includes variants of the aforementioned polypeptides, that is polypeptides that vary from the referents by conservative amino acid substitutions, whereby a residue is substituted by another with like characteristics. Typical such substitutions are among Ala, Val, Leu and Ile; among Ser and Thr; among the acidic residues Asp and Glu; among Asn and Gln; and among the basic residues Lys and Arg; or aromatic residues Phe and Tyr.\nPolypeptides of the present invention can be prepared in any suitable manner. Such polypeptides include isolated naturally occurring polypeptides, recombinantly produced polypeptides, synthetically produced polypeptides, or polypeptides produced by a combination of these methods. Means for preparing such polypeptides are well understood in the art.\nIt is most preferred that a polypeptide of the invention is derived from Moraxella catorrholis, however, it may preferably be obtained from other organisms of the same taxonomic genus. A polypeptide of the invention may also be obtained, for example, from organisms of the same taxonomic family or order.\nPolynucleotides\nIt is an object of the invention to provide polynucleotides that encode BASB023 polypeptides, particularly polynucleotides that encode the polypeptide herein designated BASB023.\nIn a particularly preferred embodiment of the invention the polynucleotide comprises a region encoding BASB023 polypeptides comprising a sequence set out in SEQ ID NO:1 or 3 which includes a full length gene, or a variant thereof.\nThe BASB023 polynucleotides provided in SEQ ID NO:1 or 3 are the BASB023 polynucleotides from Moraxella catarrhalis strains Mc293 1 (ATCC 43617).\nAs a further aspect of the invention there are provided isolated nucleic acid molecules encoding and/or expressing BASB023 polypeptides and polynucleotides, particularly Moraxella catarrhalis BASB023 polypeptides and polynucleotides, including, for example, unprocessed RNAs, ribozyme RNAs, mRNAs, cDNAs, genomic DNAs, B- and Z-DNAs. Further embodiments of the invention include biologically, diagnostically, prophylactically, clinically or therapeutically useful polynucleotides and polypeptides, and variants thereof, and compositions comprising the same.\nAnother aspect of the invention relates to isolated polynucleotides, including at least one full length gene, that encodes a BASB023 polypeptide having a deduced amino acid sequence of SEQ ID NO:2 or 4 and polynucleotides closely related thereto and variants thereof.\nIn another particularly preferred embodiment of the invention there is a BASB023 polypeptide from Moraxella catarrhalis comprising or consisting of an amino acid sequence of SEQ ID NO:2 or 4 or a variant thereof.\nUsing the information provided herein, such as a polynucleotide sequence set out in SEQ ID NO:1 or 3, a polynucleotide of the invention encoding BASB023 polypeptide may be obtained using standard cloning and screening methods, such as those for cloning and sequencing chromosomal DNA fragments from bacteria using Moraxella catarrhalis Catlin cells as starting material, followed by obtaining a full length clone. For example, to obtain a polynucleotide sequence of the invention, such as a polynucleotide sequence given in SEQ ID NO:1 or 3, typically a library of clones of chromosomal DNA of Moraxella catarrhalis Catlin in E.coli or some other suitable host is probed with a radiolabeled oligonucleotide, preferably a 17-mer or longer, derived from a partial sequence. Clones carrying DNA identical to that of the probe can then be distinguished using stringent hybridization conditions. By sequencing the individual clones thus identified by hybridization with sequencing primers designed from the original polypeptide or polynucleotide sequence it is then possible to extend the polynucleotide sequence in both directions to determine a full length gene sequence. Conveniently, such sequencing is performed, for example, using denatured double stranded DNA prepared from a plasmid clone. Suitable techniques are described by Maniatis, T., Fritsch, E. F. and Sambrook et al., MOLECULAR CLONING, A LABORATORY MANUAL, 2nd Ed.; Cold Spring Harbor Laboratory Press, Cold Spring Harbor, N.Y. (1989). (see in particular Screening By Hybridization 1.90 and Sequencing Denatured Double-Stranded DNA Templates 13.70). Direct genomic DNA sequencing may also be performed to obtain a full length gene sequence. Illustrative of the invention, each polynucleotide set out in SEQ ID NO:1 or 3 was discovered in a DNA library derived from Moraxella catarrhalis. \nMoreover, each DNA sequence set out in SEQ ID NO:1 or 3 contains an open reading frame encoding a protein having about the number of amino acid residues set forth in SEQ ID NO:2 or 4 with a deduced molecular weight that can be calculated using amino acid residue molecular weight values well known to those skilled in the art.\nThe polynucleotide of SEQ ID NO:1, between the start codon at nucleotide number 1 and the stop codon which begins at nucleotide number 808 of SEQ ID NO: 1. encodes the polypeptide of SEQ ID NO:2.\nThe polynucleotide of SEQ ID NO:3, between the start codon at nucleotide number 1 and the stop codon which begins at nucleotide number 808 of SEQ ID NO:3, encodes the polypeptide of SEQ ID NO:4.\nIn a further aspect, the present invention provides for an isolated polynucleotide comprising or consisting of:\n(a) a polynucleotide sequence which has at least 85% identity, preferably at least 90% identity, more preferably at least 95% identity, even more preferably at least 97-99% or exact identity to SEQ ID NO:1 or 3 over the entire length of SEQ ID NO:1 or 3 respectively; or\n(b) a polynucleotide sequence encoding a polypeptide which has at least 85% identity, preferably at least 90% identity, more preferably at least 95% identity, even more preferably at least 97-99% or 100% exact, to the amino acid sequence of SEQ ID NO:2 or 4, over the entire length of SEQ ID NO:2 or 4 respectively.\nA polynucleotide encoding a polypeptide of the present invention, including homologs and orthologs from species other than Moraxella catarrhalis, may be obtained by a process which comprises the steps of screening an appropriate library under stringent hybridization conditions (for example, using a temperature in the range of 45-65xc2x0 C. and an SDS concentration from 0.1-1%) with a labeled or detectable probe consisting of or comprising the sequence of SEQ ID NO:1 or 3 or a fragment thereof; and isolating a full-length gene and/or genomic clones containing said polynucleotide sequence.\nThe invention provides a polynucleotide sequence identical over its entire length to a coding sequence (open reading frame) in SEQ ID NO:1 or 3. Also provided by the invention is a coding sequence for a mature polypeptide or a fragment thereof, by itself as well as a coding sequence for a mature polypeptide or a fragment in reading frame with another coding sequence, such as a sequence encoding a leader or secretory sequence, a pre-, or pro- or prepro-protein sequence. The polynucleotide of the invention may also contain at least one non-coding sequence, including for example, but not limited to at least one noncoding 5xe2x80x2 and 3xe2x80x2 sequence, such as the transcribed but non-translated sequences, termination signals (such as rho-dependent and rho-independent termination signals), ribosome binding sites, Kozak sequences, sequences that stabilize mRNA, introns, and polyadenylation signals. The polynucleotide sequence may also comprise additional coding sequence encoding additional amino acids. For example, a marker sequence that facilitates purification of the fused polypeptide can be encoded. In certain embodiments of the invention, the marker sequence is a hexa-histidine peptide, as provided in the pQE vector (Qiagen, Inc.) and described in Gentz et al., Proc. Natl. Acad. Sci., USA 86: 821-824 (1989), or an HA peptide tag (Wilson et al., Cell 37: 767 (1984), both of which may be useful in purifying polypeptide sequence fused to them. Polynucleotides of the invention also include, but are not limited to, polynucleotides comprising a structural gene and its naturally associated sequences that control gene expression.\nThe nucleotide sequence encoding FASB023 polypeptide of SEQ ID NO:2 or 4 may be identical to the polypeptide encoding sequence contained in nucleotides 1 to 807 of SEQ ID NO:1 or 3 respectively. Alternatively it may be a sequence, which as a result of the redundancy (degeneracy) of the genetic code, also encodes the polypeptide of SEQ ID NO:2 or 4.\nThe term xe2x80x9cpolynucleotide encoding a polypeptidexe2x80x9d as used herein encompasses polynucleotides that include a sequence encoding a polypeptide of the invention, particularly a bacterial polypeptide and more particularly a polypeptide of the Moraxella catarrhalis BASB023 having an amino acid sequence set out in SEQ ID NO:2 or 4. The term also encompasses polynucleotides that include a single continuous region or discontinuous regions encoding the polypeptide (for example, polynucleotides interrupted by integrated phage, an integrated insertion sequence, an integrated vector sequence, an integrated transposon sequence, or due to RNA editing or genomic DNA reorganization) together with additional regions, that also may contain coding and/or non-coding sequences.\nThe invention further relates to variants of the polynucleotides described herein that encode variants of a polypeptide having a deduced amino acid sequence of SEQ ID NO:2 or 4. Fragments of polynucleotides of the invention may be used, for example, to synthesize full-length polynucleotides of the invention.\nFurther particularly preferred embodiments are polynucleotides encoding BASB023 variants, that have the amino acid sequence of BASB023 polypeptide of SEQ ID NO:2 or 4 in which several, a few, 5 to 10, 1 to 5, 1 to 3, 2, 1 or no amino acid residues are substituted, modified, deleted and/or added, in any combination. Especially preferred among these are silent substitutions, additions and deletions, that do not alter the properties and activities of BASB023 polypeptide.\nFurther preferred embodiments of the invention are polynucleotides trait are at least 85% identical over their entire length to a polynucleotide encoding BASB023 polypeptide having an amino acid sequence set out in SEQ ID NO:2 or 4, and polynucleotides that are complementary to such polynucleotides. Alternatively, most highly preferred are polynucleotides that comprise a region that is at least 90% identical over its entire length to a polynucleotide encoding BASB023 polypeptide and polynucleotides complementary thereto. In this regard, polynucleotides at least 95% identical over their entire length to the same are particularly preferred. Furthermore, those with at least 97% are highly preferred among those with at least 95%, and among these those with at least 98% and at least 99% are particularly highly preferred, with at least 99% being the more preferred.\nPreferred embodiments are polynucleotides encoding polypeptides that retain substantially the same biological function or activity as the mature polypeptide encoded by a DNA of SEQ ID NO:1 or 3.\nIn accordance with certain preferred embodiments of this invention there are provided polynucleotides that hybridize, particularly under stringent conditions, to BASB023 polynucleotide sequences, such as those polynucleotides in SEQ ID NO:1 or 3.\nThe invention further relates to polynucleotides that hybridize to the polynucleotide sequences provided herein. In this regard, the invention especially relates to polynucleotides that hybridize under stringent conditions to the polynucleotides described herein. As herein used, the terms xe2x80x9cstringent conditionsxe2x80x9d and xe2x80x9cstringent hybridization conditionsxe2x80x9d mean hybridization occurring only if there is at least 95% and preferably at least 97% identity between the sequences. A specific example of stringent hybridization conditions is overnight incubation at 42xc2x0 C. in a solution comprising: 50%/o formamide, 5xc3x97SSC (150 mM NaCl, 15 mM trisodium citrate), 50 mM sodium phosphate (pH7.6), 5xc3x97Denhardt\"\"s solution, 10% dextran sulfate, and 20 micrograms/ml of denatured, sheared salmon sperm DNA, followed by washing the hybridization support in 0.1xc3x97SSC at about 65xc2x0 C. Hybridization and wash conditions are well known and exemplified in Sambrook, et al., Molecular Cloning: A Laboratory Manual, Second Edition, Cold Spring Harbor, N.Y. (1989), particularly Chapter 11 therein. Solution hybridization may also be used with the polynucleotide sequences provided by the invention.\nThe invention also provides a polynucleotide consisting of or comprising a polynucleotide sequence obtained by screening an appropriate library containing the complete gene for a polynucleotide sequence set forth in SEQ ID NO:1 or 3 under stringent hybridization conditions with a probe having the sequence of said polynucleotide sequence set forth in SEQ ID NO:1 or 3 or a fragment thereof; and isolating said polynucleotide sequence. Fragments useful for obtaining such a polynucleotide include, for example, probes and primers fully described elsewhere herein.\nAs discussed elsewhere herein regarding polynucleotide assays of the invention, for instance, the polynucleotides of the invention, may be used as a hybridization probe for RNA, cDNA and genomic DNA to isolate full-length cDNAs and genomic clones encoding BASB023 and to isolate cDNA and genomic clones of other genes that have a high identity, particularly high sequence identity, to the BASB023 gene. Such probes generally will comprise at least 15 nucleotide residues or base pairs. Preferably, such probes will have at least 30 nucleotide residues or base pairs and may have at least 50 nucleotide residues or base pairs. Particularly preferred probes will have at least 20 nucleotide residues or base pairs and will have less than 30 nucleotide residues or base pairs.\nA coding region of a BASB023 gene may be isolated by screening using a DNA sequence provided in SEQ ID NO:1 or 3 to synthesize an oligonucleotide probe. A labeled oligonucleotide having a sequence complementary to that of a gene of the invention is then used to screen a library of cDNA, genomic DNA or mRNA to determine which members of the library the probe hybridizes to.\nThere are several methods available and well known to those skilled in the art to obtain full-length DNAs, or extend short DNAs, for example those based on the method of Rapid Amplification of cDNA ends (RACE) (see, for example, Frohman, et al., PNAS USA 85: 8998-9002. 1988). Recent modifications of the technique, exemplified by the Marathon(trademark) technology (Clontech Laboratories Inc.) for example, have significantly simplified the search for longer cDNAs. In the Marathon(trademark) technology, cDNAs have been prepared from mRNA extracted from a chosen tissue and an xe2x80x98adaptorxe2x80x99 sequence ligated onto each end. Nucleic acid amplification (PCR) is then carried out to amplify the xe2x80x9cmissingxe2x80x9d 5xe2x80x2 end of the DNA using a combination of gene specific and adaptor specific oligonucleotide primers. The PCR reaction is then repeated using xe2x80x9cnestedxe2x80x9d primers, that is, primers designed to anneal within the amplified product (typically an adaptor specific primer that anneals further 3xe2x80x2 in the adaptor sequence and a gene specific primer that anneals further 5xe2x80x2 in the selected gene sequence). The products of this reaction can then be analyzed by DNA sequencing and a full-length DNA constructed either by joining the product directly to the existing DNA to give a complete sequence, or carrying out a separate full-length PCR using the new sequence information for the design of the 5xe2x80x2 primer.\nThe polynucleotides and polypeptides of the invention may be employed, for example, as research reagents and materials for discovery of treatments of and diagnostics for diseases, particularly human diseases, as further discussed herein relating to polynucleotide assays.\nThe polynucleotides of the invention that are oligonucleotides derived from a sequence of SEQ ID NOS:1-4 may be used in the processes herein as described, but preferably for PCR, to determine whether or not the polynucleotides identified herein in whole or in part are transcribed in bacteria in infected tissue. It is recognized that such sequences will also have utility in diagnosis of the stage of infection and type of infection the pathogen has attained.\nThe invention also provides polynucleotides that encode a polypeptide that is the mature protein plus additional amino or carboxyl-terminal amino acids, or amino acids interior to the mature polypeptide (when the mature form has more than one polypeptide chain, for instance). Such sequences may play a role in processing of a protein from precursor to a mature form, may allow protein transport may lengthen or shorten protein half-life or may facilitate manipulation of a protein for assay or production, among other things. As generally is the case in vivo, the additional amino acids may be processed away from the mature protein by cellular enzymes.\nFor each and every polynucleotide of the invention there is provided a polynucleotide complementary to it. It is preferred that these complementary polynucleotides are fully complementary to each polynucleotide with which they are complementary.\nA precursor protein, having a mature form of the polypeptide fused to one or more prosequences may be an inactive form of the polypeptide. When prosequences are removed such inactive precursors generally are activated. Some or all of the prosequences may be removed before activation. Generally, such precursors are called proproteins.\nIn addition to the standard A, G, C, T/U representations for nucleotides, the term xe2x80x9cNxe2x80x9d may also be used in describing certain polynucleotides of the invention. xe2x80x9cNxe2x80x9d means that any of the four DNA or RNA nucleotides may appear at such a designated position in the DNA or RNA sequence, except it is preferred that N is not a nucleic acid that when taken in combination with adjacent nucleotide positions, when read in the correct reading frame, would have the effect of generating a premature termination codon in such reading frame.\nIn sum, a polynucleotide of the invention may encode a mature protein, a mature protein plus a leader sequence (which may be referred to as a preprotein), a precursor of a mature protein having one or more prosequences that are not the leader sequences of a preprotein, or a preproprotein, which is a precursor to a proprotein, having a leader sequence and one or more prosequences, which generally are removed during processing steps that produce active and mature forms of the polypeptide.\nIn accordance with an aspect of the invention, there is provided the use of a polynucleotide of the invention for therapeutic or prophylactic purposes, in particular genetic immunization.\nThe use of a polynucleotide of the invention in genetic immunization will preferably employ a suitable delivery method such as direct injection of plasmid DNA into muscles (Wolff et al., Hum Mol Genet (1992) 1:363, Manthorpe et al., Hum. Gene Ther. (1983) 4:419), delivery of DNA complexed with specific protein carriers (Wu et al., J Biol Chem. (1989) 264: 16985), coprecipitation of DNA with calcium phosphate (Benvenisty and Reshef, PNAS USA, (1986) 83: 9551), encapsulation of DNA in various forms of liposomes (Kaneda et al., Science (1989) 243: 375), particle bombardment (Tang et al., Nature (1992) 356:152, Eisenbraun et al., DNA Cell Biol (1993) 12: 791) and in vivo infection using cloned retroviral vectors (Seeger et al., PNAS USA (1984) 81: 5849).\nVectors, Host Cells, Expression Systems\nThe invention also relates to vectors that comprise a polynucleotide or polynucleotides of the invention, host cells that are genetically engineered with vectors of the invention and the production of polypeptides of the invention by recombinant techniques. Cell-free translation systems can also be employed to produce such proteins using RNAs derived from the DNA constructs of the invention.\nRecombinant polypeptides of the present invention may be prepared by processes well known in those skilled in the art from genetically engineered host cells comprising expression systems. Accordingly, in a further aspect, the present invention relates to expression systems that comprise a polynucleotide or polynucleotides of the present invention, to host cells which are genetically engineered with such expression systems, and to the production of polypeptides of the invention by recombinant techniques.\nFor recombinant production of the polypeptides of the invention, host cells can be genetically engineered to incorporate expression systems or portions thereof or polynucleotides of the invention. Introduction of a polynucleotide into the host cell can be effected by methods described in many standard laboratory manuals, such as Davis, et al., BASIC METHODS IN MOLECULAR BIOLOGY, (1986) and Sambrook, et al., MOLECULAR CLONING: A LABORATORY MANUAL, 2nd Ed., Cold Spring Harbor Laboratory Press, Cold Spring Harbor, N.Y. (1989), such as, calcium phosphate transfection, DEAE-dextran mediated transfection transvection, microinjection, cationic lipid-mediated transfection, electroporation, transduction, scrape loading, ballistic introduction and infection.\nRepresentative examples of appropriate hosts include bacterial cells, such as cells of streptococci, staphylococci, enterococci, E. coli, streptomyces, cyanobacteria, Bacillus subtilis, Neisseria meningitidis and Moraxella catarrhalis; fungal cells, such as cells of a yeast, Kluveromyces, Saccharomyces, a basidiomycete, Candida albicans and Aspergillus; insect cells such as cells of Drosophila S2 and Spodoptera Sf9; animal cells such as CHO, COS, HeLa, C127, 3T3, BHK, 293, CV-1 and Bowes melanoma cells; and plant cells, such as cells of a gymnosperm or angiosperm.\nA great variety of expression systems can be used to produce the polypeptides of the invention. Such vectors include, among others, chromosomal-, episomal- and virus-derived vectors, for example, vectors derived from bacterial plasmids, from bacteriophage, from transposons, from yeast episomes, from insertion elements, from yeast chromosomal elements, from viruses such as baculoviruses, papova viruses, such as SV40, vaccinia viruses, adenoviruses, fowl pox viruses, pseudorabies viruses, picomaviruses, retroviruses, and alphaviruses and vectors derived from combinations thereof, such as those derived from plasmid and bactcriophage genetic elements, such as cosmids and phagemids. The expression system constructs may contain control regions that regulate as well as engender expression. Generally, any system or vector suitable to maintain, propagate or express polynucleotides and/or to express a polypeptide in a host may be used for expression in this regard. The appropriate DNA sequence may be inserted into the expression system by any of a variety of well-known and routine techniques, such as, for example, those set forth in Sambrook et al., MOLECULAR CLONING, A LABORATORY MANUAL, (supra).\nIn recombinant expression systems in eukaryotes, for secretion of a translated protein into the lumen of the endoplasmic reticulum, into the periplasmic space or into the extracellular environment, appropriate secretion signals may be incorporated into the expressed polypeptide. These signals may be endogenous to the polypeptide or they may be heterologous signals.\nPolypeptides of the present invention can be recovered and purified from recombinant cell cultures by well-known methods including ammonium sulfate or ethanol precipitation, acid extraction, anion or cation exchange chromatography, phosphocellulose chromatography, hydrophobic interaction chromatography, affinity chromatography, hydroxylapatite chromatography and lectin chromatography. Most preferably, ion metal affinity chromatography (GMAC) is employed for purification. Well known techniques for refolding proteins may be employed to regenerate active conformation when the polypeptide is denatured during intracellular synthesis, isolation and or purification.\nThe expression system may also be a recombinant live microorganism, such as a virus or bacterium. The gene of interest can be inserted into the genome of a live recombinant virus or bacterium. Inoculation and in vivo infection with this live vector will lead to in vivo expression of the antigen and induction of immune responses. Viruses and bacteria used for this purpose are for instance: poxviruses (e.g.; vaccinia, fowlpox, canarypox), alphaviruses (Sindbis virus, Semliki Forest Virus, Venezuelian Equine Encephalitis Virus), adenoviruses, adeno-associated virus, picornaviruses (poliovirus, rhinovirus). herpesviruses (varicella zoster virus, etc), Listeria, Salmonella, Shigella, BCG. These viruses and bacteria can be virulent, or attenuated in various ways in order to obtain live vaccines. Such live vaccines also form part of the invention.\nDiagnostic, Prognostic, Serotyping and Mutation Assays\nThis invention is also related to the use of BASB023 polynucleotides and polypeptides of the invention for use as diagnostic reagents. Detection of BASB023 polynucleotides and/or polypeptides in a eukaryote, particularly a mammal, and especially a human, will provide a diagnostic method for diagnosis of disease, staging of disease or response of an infectious organism to drugs. Eukaryotes, particularly mammals, and especially humans, particularly those infected or suspected to be infected with an organism comprising the BASB023 gene or protein, may be detected at the nucleic acid or amino acid level by a variety of well known techniques as cell as by methods provided herein.\nPolypeptides and polynucleotides for prognosis, diagnosis or other analysis may be obtained from a putatively infected and/or infected individual\"\"s bodily materials. Polynucleotides from any of these sources, particularly DNA or RNA, may be used directly for detection or may be amplified enzymatically by using PCR or any other amplification technique prior to analysis. RNA, particularly mRNA, cDNA and genomic DNA may also be used in the same ways. Using amplification, characterization of the species and strain of infectious or resident organism present in an individual, may be made by an analysis of the genotype of a selected polynucleotide of the organism. Deletions and insertions can be detected by a change in size of the amplified product in comparison to a genotype of a reference sequence selected from a related organism, preferably a different species of the same genus or a different strain of the same species. Point mutations can be identified by hybridizing amplified DNA to labeled BASB023 polynucleotide sequences. Perfectly or significantly matched sequences can be distinguished from imperfectly or more significantly mismatched duplexes by DNase or RNase digestion, for DNA or RNA respectively, or by detecting differences in melting temperatures or renaturation kinetics. Polynucleotide sequence differences may also be detected by alterations in the electrophoretic mobility of polynucleotide fragments in gels as compared to a reference sequence. This may be carried out with or without denaturing agents. Polynucleotide differences may also be detected by direct DNA or RNA sequencing. See, for example, Myers et al., Science, 230: 1242 (1985). Sequence changes at specific locations also may be revealed by nuclease protection assays, such as RNase. V1 and S1 protection assay or a chemical cleavage method. See, for example, Cotton et al., Proc. Natl. Acad. Sci., USA, 85: 4397-4401 (1985).\nIn another embodiment, an array of oligonucleotides probes comprising BASB023 nucleotide sequence or fragments thereof can be constructed to conduct efficient screening of, for example, genetic mutations, serotype, taxonomic classification or identification. Array technology methods are well known and have general applicability and can be used to address a variety of questions in molecular genetics including gene expression, genetic linkage, and genetic variability (see, for example, Chee et al., Science, 274: 610 (1996)).\nThus in another aspect, the present invention relates to a diagnostic kit which comprises:\n(a) a polynucleotide of the present invention, preferably the nucleotide sequence of SEQ ID NO:1 or 3, or a fragment thereof;\n(b) a nucleotide sequence complementary to that of (a);\n(c) a polypeptide of the present invention, preferably the polypeptide of SEQ ID NO:2 or 4 or a fragment thereof; or\n(d) an antibody to a polypeptide of the present invention, preferably to the polypeptide of SEQ ID NO:2 or 4.\nIt will be appreciated that in any such kit, (a), (b), (c) or (d) may comprise a substantial component. Such a kit will be of use in diagnosing a disease or susceptibility to a Disease, among others.\nThis invention also relates to the use of polynucleotides of the present invention as diagnostic reagents. Detection of a mutated form of a polynucleotide of the invention, preferable, SEQ ID NO:1 or 3, which is associated with a disease or pathogenicity will provide a diagnostic tool that can add to, or define, a diagnosis of a disease, a prognosis of a course of disease, a determination of a stage of disease, or a susceptibility to a disease, which results from under-expression, over-expression or altered expression of the polynucleotide. Organisms, particularly infectious organisms, carrying mutations in such polynucleotide may be detected at the polynucleotide level by a variety of techniques, such as those described elsewhere herein.\nCells from an organism carrying mutations or polymorphisms (allelic variations) in a polynucleotide and/or polypeptide of the invention may also be detected at the polynucleotide or polypeptide level by a variety of techniques, to allow for serotyping, for example. For example, RT-PCR can be used to detect mutations in the RNA. It is particularly preferred to use RT-PCR in con unction with automated detection systems, such as, for example, GeneScan. RNA, cDNA or genomic DNA may also be used for the same purpose, PCR. As an example, PCR primers complementary to a polynucleotide encoding BASB023 polypeptide can be used to identify and analyze mutations.\nThe invention further provides primers with 1, 2, 3 or 4 nucleotides removed from the 5xe2x80x2 and/or the 3xe2x80x2 end. These primers may be used for, among other things, amplifying BASB023 DNA and/or RNA isolated from a sample derived from an individual, such as a bodily material. The primers may be used to amplify a polynucleotide isolated from an infected individual, such that the polynucleotide may then be subject to various techniques for elucidation of the polynucleotide sequence. In this way, mutations in the polynucleotide sequence may be detected and used to diagnose and/or prognose the infection or its stage or course, or to serotype and/or classify the infectious agent.\nThe invention further provides a process for diagnosing, disease, preferably bacterial infections, more preferably infections caused by Moraxella catarrhalis, comprising, determining from a sample derived from an individual, such as a bodily material, an increased level of expression of polynucleotide having a sequence of SEQ ID NO:1 or 3. Increased or decreased expression of a BASB023 polynucleotide can be measured using any on of the methods well known in the art for the quantitation of polynucleotides, such as, for example, amplification. PCR, RT-PCR, RNase protection, Northern blotting, spectrometry and other hybridization methods.\nIn addition, a diagnostic assay in accordance with the invention for detecting over-expression of BASB023 polypeptide compared to normal control tissue samples may be used to detect the presence of an infection, for example. Assay techniques that can be used to determine levels of a BASB023 polypeptide, in a sample derived from a host, such as a bodily material, are well-known to those of skill in the art. Such assay methods include radioimmunoassays, competitive-binding assays, Western Blot analysis, antibody sandwich assays, antibody detection and ELISA assays.\nThe polynucleotides of the invention may be used as components of polynucleotide arrays, preferably high density arrays or grids. These high density arrays are particularly useful for diagnostic and prognostic purposes. For example, a set of spots each comprising a different gene, and further comprising a polynucleotide or polynucleotides of the invention, may be used for probing, such as using hybridization or nucleic acid amplification, using a probes obtained or derived from a bodily sample, to determine the presence of a particular polynucleotide sequence or related sequence in an individual. Such a presence may indicate the presence of a pathogen, particularly Moraxella catarrhalis, and may be useful in diagnosing and/or prognosing disease or a course of disease. A grid comprising a number of variants of the polynucleotide sequence of SEQ ID NO:1 or 3 are preferred. Also preferred is a comprising a number of variants of a polynucleotide sequence encoding the polypeptide sequence of SEQ ID NO:2 or 4.\nAntibodies\nThe polypeptides and polynucleotides of the invention or variants thereof, or cells expressing the same can be used as immunogens to produce antibodies immunospecific for such polypeptides or polynucleotides respectively.\nIn certain preferred embodiments of the invention there are provided antibodies against BASB023 polypeptides or polynucleotides.\nAntibodies generated against the polypeptides or polynucleotides of the invention can be obtained by administering the polypeptides and/or polynucleotides of the invention, or epitope-bearing fragments of either or both, analogues of either or both, or cells expressing either or both, to an animal, preferably a nonhuman, using routine protocols. For preparation of monoclonal antibodies, any technique known in the art that provides antibodies produced by continuous cell line cultures can be used. Examples include various techniques, such as those in Kohler, G. and Milstein, C., Nature 256: 495-497 (1975); Kozbor et al., Immunology Today 4: 72 (1983); Cole et al., pg. 77-96 in MONOCLONAL ANTIBODIES AND CANCER THERAPY, Alan R. Liss, Inc. (1985).\nTechniques for the production of single chain antibodies (U.S. Pat. No. 4,946,778) can be adapted to produce single chain antibodies to polypeptides or polynucleotides of this invention. Also, transgenic mice, or other organisms or animals, such as other mammals, may be used to express humanized antibodies immunospecific to the polypeptides or polynucleotides of the invention.\nAlternatively, phage display technology may be utilized to select antibody genes with binding activities towards a polypeptide of the invention either from repertoires of PCR amplified v-genes of lymphocytes from humans screened for possessing anti-BASB023 or from naive libraries (McCafferty, et al., (1990), Nature 348, 552-554Marks. et al., (1992) Biotechnology 10. 779-783). The affinity of these antibodies can also be improved by, for example, chain shuffling (Clackson et al., (1991) Nature 352: 628).\nThe above-described antibodies may be employed to isolate or to identify clones expressing the polypeptides or polynucleotides of the invention to purify the polypeptides or polynucleotides by, for example, affinity chromatography.\nThus, among others, antibodies against BASB023-polypeptide or BASB023-polynucleotide may be employed to treat infections, particularly bacterial infections.\nPolypeptide variants include antigenically, epitopically or immunologically equivalent variants form a particular aspect of this invention.\nPreferably, the antibody or variant thereof is modified to make it less immunogenic in the individual. For example, if the individual is human the antibody may most preferably be xe2x80x9chumanized,xe2x80x9d where the complimentarity determining region or regions of the hybridoma-derived antibody has been transplanted into a human monoclonal antibody, for example as described in Jones et al. (1986), Nature 321, 522-525 or Tempest et al., (1991) Biotechnology 9, 266-273.\nAntagonists and Agonistsxe2x80x94Assays and Molecules\nPolypeptides and polynucleotides of the invention may also be used to assess the binding of small molecule substrates and ligands in, for example, cells, cell-free preparations, chemical libraries, and natural product mixtures. These substrates and ligands may be natural substrates and ligands or may be structural or functional mimetics. See, e.g., Coligan et al., Current Protocols in Immunology I(2): Chapter 5 (1991).\nThe screening methods may simply measure the binding of a candidate compound to the polypeptide or polynucleotide, or to cells or membranes bearing the polypeptide or polynucleotide, or a fusion protein of the polypeptide by means of a label directly or indirectly associated with the candidate compound. Alternatively, the screening method may involve competition with a labeled competitor. Further, these screening methods may test whether the candidate compound results in a signal generated by activation or inhibition of the polypeptide or polynucleotide, using detection systems appropriate to the cells comprising the polypeptide or polynucleotide. Inhibitors of activation are generally assayed in the presence of a known agonist and the effect on activation by the agonist by the presence of the candidate compound is observed. Constitutively active polypeptide and/or constitutively expressed polypeptides and polynucleotides may be employed in screening methods for inverse agonists or inhibitors, in the absence of an agonist or inhibitor, by testing whether the candidate compound results in inhibition of activation of the polypeptide or polynucleotide, as the case may be. Further, the screening methods may simply comprise the steps of mixing a candidate compound with a solution containing a polypeptide or polynucleotide of the present invention, to form a mixture, measuring BASB023 polypeptide and/or polynucleotide activity in the mixture, and comparing the BASB023 polypeptide and/or polynucleotide activity of the mixture to a standard. Fusion proteins, such as those made from Fc portion and BASB023 polypeptide, as hereinbefore described, can also be used for high-throughput screening assays to identify antagonists of the polypeptide of the present invention, as well as of phylogenetically and and/or functionally related polypeptides (see D. Bennett et al., J Mol Recognition, 8:52-58 (1995); and K. Johanson et al., J Biol Chem, 270(16):9459-9471 (1995)).\nThe polynucleotides, polypeptides and antibodies that bind to and/or interact with a polypeptide of the present invention may also be used to configure screening methods for detecting the effect of added compounds on the production of mRNA and/or polypeptide in cells. For example, an ELISA assay may be constructed for measuring secreted or cell associated levels of polypeptide using monoclonal and polyclonal antibodies by standard methods known in the art. This can be used to discover agents which may inhibit or enhance the production of polypeptide (also called antagonist or agonist, respectively) from suitably manipulated cells or tissues.\nThe invention also provides a method of screening compounds to identify those which enhance (agonist) or block (antagonist) the action of BASB023 polypeptides or polynucleotides, particularly those compounds that are baeteriostatic and/or bactericidal. The method of screening may involve high-throughput techniques. For example, to screen for agonists or antagonists, a synthetic reaction mix, a cellular compartment, such as a membrane, cell envelope or cell wall, or a preparation of any thereof, comprising BASB023 polypeptide and a labeled substrate or ligand of such polypeptide is incubated in the absence or the presence of a candidate molecule that may be a BASB023 agonist or antagonist. The ability of the candidate molecule to agonize or antagonize the BASB023 polypeptide is reflected in decreased binding of the labeled ligand or decreased production of product from such substrate. Molecules that bind gratuitously, ie., without inducing the effects of BASB023 polypeptide are most likely to be good antagonists. Molecules that bind well and, as the case may be, increase the rate of product production from substrate, increase signal transduction, or increase chemical channel activity are agonists. Detection of the rate or level of, as the case may be, production of product from substrate, signal transduction, or chemical channel activity may be enhanced by using a reporter system. Reporter systems that may be useful in this regard include but are not limited to calorimetric, labeled substrate converted into product, a reporter gene that is responsive to changes in BASB023 polynucleotide or polypeptide activity, and binding assays known in the art.\nAnother example of an assay for BASB023 agonists is a competitive assay that combines BASB023 and a potential agonist with BASB023-binding molecules, recombinant BASB023 binding molecules, natural substrates or ligands, or substrate or ligand mimetics, under appropriate conditions for a competitive inhibition assay. BASB023 can be labeled, such as by radioactivity or a calorimetric compound, such that the number of BASB023 molecules bound to a binding molecule or converted to product can be determined accurately to assess the effectiveness of the potential antagonist.\nPotential antagonists include, among others, small organic molecules, peptides, polypeptides and antibodies that bind to a polynucleotide and/or polypeptide of the invention and thereby inhibit or extinguish its activity or expression. Potential antagonists also may be small organic molecules, a peptide, a polypeptide such as a closely related protein or antibody that binds the same sites on a binding molecule, such as a binding molecule, without inducing BASB023-induced activities, thereby preventing the action or expression of BASB023 polypeptides and/or polynucleotides by excluding BASB023 polypeptides and/or polynucleotides from binding.\nPotential antagonists include a small molecule that binds to and occupies the binding site of the polypeptide thereby preventing binding to cellular binding molecules, such that normal biological activity is prevented. Examples of small molecules include but are not limited to small organic molecules, peptides or peptide-like molecules. Other potential antagonists include antisense molecules (see Okano, J. Neurochem. 56: 560 (1991); OLIGODEOXYNUCLEOTIDES AS ANTISENSE INHIBITORS OF GENE EXPRESSION, CRC Press, Boca Raton, Fla. (1988), for a description of these molecules). Preferred potential antagonists include compounds related to and variants of BASB023.\nIn a further aspect, the present invention relates to genetically engineered soluble fusion proteins comprising a polypeptide of the present invention, or a fragment thereof, and various portions of the constant regions of heavy or light chains of immunoglobulins of various subclasses (IgG, IgM, IgA, IgE). Preferred as an immunoglobulin is the constant part of the heavy chain of human IgG, particularly IgG1, where fusion takes place at the hinge region. In a particular embodiment, the Fc plart can be removed simply by incorporation of a cleavage sequence which can be cleaved with blood clotting factor Xa. Furthermore, this invention relates to processes for the preparation of these fusion proteins by genetic engineering, and to the use thereof for drug screening, diagnosis and therapy. A further aspect of the invention also relates to polynucleotides encoding such fusion proteins. Examples of fusion protein technology can be found in International Patent Application Nos. WO94/29458 and WO94/22914.\nEach of the polynucleotide sequences provided herein may be used in the discovery and development of antibacterial compounds. The encoded protein, upon expression, can be used as a target for the screening of antibacterial drugs. Additionally, the polynucleotide sequences encoding the amino terminal regions of the encoded protein or Shine-Delgarno or other translation facilitating sequences of the respective mRNA can be used to construct antisense sequences to control the expression of the coding sequence of interest.\nThe invention also provides the use of the polypeptide, polynucleotide, agonist or antagonist of the invention to interfere with the initial physical interaction between a pathogen or pathogens and a eukaryotic, preferably mammalian, host responsible for sequelae of infection. In particular, the molecules of the invention may be used: in the prevention of adhesion of bacteria, in particular gram positive and/or gram negative bacteria, to eukaryotic, preferably mammalian, extracellular matrix proteins on in-dwelling devices or to extracellular matrix proteins in wounds; to block bacterial adhesion between eukaryotic, preferably mammalian, extracellular matrix proteins and bacterial BASB023 proteins that mediate tissue damage and/or; to block the normal progression of pathogenesis in infections initiated other than by the implantation of in-dwelling devices or by other surgical techniques.\nIn accordance with yet another aspect of the invention, there are provided BASB023 agonists and antagonists, preferably bacteristatic or bactericidal agonists and antagonists.\nThe antagonlists and agonists of the invention may be employed, for instance, to prevent, inhibit and or treat diseases.\nIn a further aspect, the present invention relates to mimotopes of the polypeptide of the invention. A mimotope is a peptide sequence, sufficiently similar to the native peptide (sequentially or structurally), which is capable of being recognised by antibodies which recognise the native peptide; or is capable of raising antibodies which recognise the native peptide when coupled to a suitable carrier.\nPeptide mimotopes may be designed for a particular purpose by addition, deletion or substitution of elected amino acids. Thus, the peptides may be modified for the purposes of ease of conjugation to a protein carrier. For example, it may be desirable for some chemical conjugation methods to include a terminal cysteine. In addition it may be desirable for peptides conjugated to a protein carrier to include a hydrophobic terminus distal from the conjugated terminus of the peptide, such that the free unconjugated end of the peptide remains associated with the surface of the carrier protein. Thereby presenting the peptide in a conformation which most closely resembles that of the peptide as found in the context of the whole native molecule. For example, the peptides may be altered to have an N-terminal cysteine and a C-terminal hydrophobic amidated tail. Alternatively, the addition or substitution of a D-stereoisomer form of one or more of the amino acids may be performed to create a beneficial derivative, for example to enhance stability of the peptide.\nAlternatively, peptide mimotopes may be identified using antibodies which are capable themselves of binding to the polypeptides of the present invention using techniques such as phage display technology (EP 0 552 267 B1). This technique, generates a large number of peptide sequences which mimic the structure of the native peptides and are, therefore, capable of binding to anti-native peptide antibodies, but may not necessarily themselves share significant sequence homology to the native polypeptide.\nVaccines\nAnother aspect of the invention relates to a method for inducing an immunological response in an individual, particularly a mammal, preferably humans, which comprises inoculating the individual with BASB023 polynucleotide and/or polypeptide, or a fragment or variant thereof, adequate to produce antibody and/ or T cell immune response to protect said individual from infection, particularly bacterial infection and most particularly Moraxella catarrhalis infection. Also provided are methods whereby such immunological response slows bacterial replication. Yet another aspect of the invention relates to a method of inducing immunological response in an individual which comprises delivering to such individual a nucleic acid vector, sequence or ribozyme to direct expression of BASB023 polynucleotide and/or polypeptide, or a fragment or a variant thereof, for expressing BASB023 polynucleotide and/or polypeptide, or a fragment or a variant thereof in vivo in order to induce an immunological response, such as, to produce antibody and/ or T cell immune response, including, for example, cytokine-producing T cells or cytotoxic T cells, to protect said individual, preferably a human, from disease, whether that disease is already established within the individual or not. One example of administering the gene is by accelerating it into the desired cells as a coating on particles or otherwise. Such nucleic acid vector may comprise DNA, RNA, a ribozyme, a modified nucleic acid, a DNA/RNA hybrid, a DNA-protein complex or an RNA-protein complex.\nA further aspect of the invention relates to an immunological composition that when introduced into an individual, preferably a human, capable of having induced within it an immunological response, induces an immunological response in such individual to a BASB023 polynucleotide and/or polypeptide encoded therefrom, wherein the composition comprises a recombinant BASB023 polynucleotide and/or polypeptide encoded therefrom and/or comprises DNA and/or RNA which encodes and expresses an antigen of said BASB023 polynucleotide, polypeptide encoded therefrom, or other polypeptide of the invention. The immunological response may be used therapeutically or prophylactically and may take the form of antibody immunity and/or cellular immunity, such as cellular immunity arising from CTL or CD4+T cells.\nA BASB023 polypeptide or a fragment thereof may be fused with co-protein or chemical moiety which may or may not by itself produce antibodies, but which is capable of stabilizing the first protein and producing a fused or modified protein which will have antigenic and/or immunogenic properties, and preferably protective properties. Thus fused recombinant protein, preferably further comprises an antigenic co-protein, such as lipoprotein D from Haemophilius influenzae, Glutathione-S-transferase (GST) or beta-galactosidase, or any other relatively large co-protein which solubilizes the protein and facilitates production and purification thereof. Moreover, the co-protein may act as an adjuvant in the sense of providing a generalized stimulation of the immune system of the organism receiving the protein. The co-protein may be attached to either the amino- or carboxy-terminus of the first protein.\nProvided by this invention are compositions, particularly vaccine compositions, and methods comprising the polypeptides and/or polynucleotides of the invention and immunostimulatory DNA sequences, such as those described in Sato, Y. et al. Science 273: 352 (1996).\nAlso, provided by this invention are methods using the described polynucleotide or particular fragments thereof, which have been shown to encode non-variable regions of bacterial cell surface proteins, in polynucleotide constructs used in such genetic immunization experiments in animal models of infection with Moraxelia catarrhalis. Such experiments will be particularly useful for identifying protein epitopes able to provoke a prophylactic or therapeutic immune response. It is believed that this approach will allow for the subsequent preparation of monoclonal antibodies of particular value, derived from the requisite organ of the animal successfully resisting or clearing infection, for the development of prophylactic agents or therapeutic treatments of bacterial infection, particularly Moraxella catarrhalis infection, in mammals, particularly humans.\nThe invention also includes a vaccine formulation which comprises an immunogenic recombinant polypeptide and/or polynucleotide of the invention together with a suitable carrier, such as a pharmaceutically acceptable carrier. Since the polypeptides and polynucleotides may be broken down in the stomach, each is preferably administered parenterally, including, for example, administration that is subcutaneous, intramuscular, intravenous, or intradermal. Formulations suitable for parenteral administration include aqueous and non-aqueous sterile injection solutions which may contain anti-oxidants, buffers, bacteriostatic compounds and solutes which render the formulation isotonic with the bodily fluid, preferably the blood, of the individual; and aqueous and non-aqueous sterile suspensions which may include suspending agents or thickening agents. The formulations may be presented in unit-dose or multi-dose containers, for example, sealed ampoules and vials and may be stored in a freeze-dried condition requiring only the addition of the sterile liquid carrier immediately prior to use.\nThe vaccine formulation of the invention may also include adjuvant systems for enhancing the immunogenicity of the formulation. Preferably the adjuvant system raises preferentially a TH1 type of response.\nAn immune response may be broadly distinguished into two extreme catagories, being a humoral or cell mediated immune responses (traditionally characterised by antibody and cellular effector mechanisms of protection respectively). These categories of response have been termed TH1-type responses (cell-mediated response), and TH2-type immune responses (humoral response).\nExtreme TH1-type immune responses may be charactcrised by the generation of antigen specific. haplotype restricted cytotoxic T lymphocytes, and natural killer cell responses.\nIn mice TH1-type responses are often characterised by the generation of antibodies of the IgG2a subtype. whilst in the human these correspond to IgG1 type antibodies. TH2-type immune responses are characterised by the generation of a broad range of immunoglobulin isotypes including in mice IgG1, IgA, and IgM.\nIt can be considered that the driving force behind the development of these two types of immune responses are cytokines. High levels of TH1-type cytokines tend to favour the induction of cell mediated immune responses to the given antigen, whilst high levels of TH2-type cytokines tend to favour the induction of humoral immune responses to the antigen.\nThe distinction of TH1 and TH2-type immune responses is not absolute. In reality an individual will support an immune response which is described as being predominantly TH1 or predominantly TH2. However, it is often convenient to consider the families of cytokines in terms of that described in murine CD4+ve T cell clones by Mosmann and Coffman (Mosmana T. R. and Coffman, R. L. (1989) TH1 and TH2 cells: different patterns of lymphokine secretion lead to different functional properties. Annual Review of Immunology, 7, p145-173). Traditionally, TH1-type responses are associated with the production of the INF-xcex3 and IL-2 cytokines by T-lymphocytes. Other cytokines often directly associated with the induction of TH1-type immune responses are not produced by T-cells. such as IL-12. In contrast, TH2- type responses are associated with the secretion of IL4, IL-5, IL-6 and IL-13.\nIt is known that certain vaccine adjuvants are particularly suited to the stimulation of either TH1 or TH2- type cytokine responses. Traditionally the best indicators of the TH1:TH2 balance of the immune response after a vaccination or infection includes direct measurement of the production of TH1 or TH2 cytokines by T lymphocytes in vitro after restimulation with antigen, and/or the measurement of the IgG1:IgG2a ratio of antigen specific antibody responses.\nThus, a TH1-type adjuvant is one which preferentially stimulates isolated T-cell populations to produce high levels of TH1-type cytokines when re-stimulated with antigen in vitro, and promotes development of both CD8+ cytotoxic T lymphocytes and antigen specific immunoglobulin responses associated with TH1-type isotype.\nAdjuvants which are capable of preferential stimulation of the TH1 cell response are described in International Patent Application No. WO 94/00153 and WO 95/17209.\n3 De-O-acylated monophosphoryl lipid A (3D-MPL) is one such adjuvant. This is known from GB 2220211 (Ribi). Chemically it is a mixture of 3 De-O-acylated monophosphoryl lipid A with 4, 5 or 6 acylated chains and is manufactured by Ribi Immunochem, Montana. A preferred form of 3 De-O-acylated monophosphoryl lipid A is disclosed in European Patent 0 689 454 B1 (SmithKline Beecham Biologicals SA).\nPreferably, the particles of 3D-MPL are small enough to be sterile filtered through a 0.22 micron membrane (European Patent number 0 689 454). 3D-MPL will be present in the range of 10 xcexcg-100 xcexcg preferably 25-50 xcexcg per dose wherein the antigen will typically be present in a range 2-50 xcexcg per dose.\nAnother preferred adjuvant comprises QS21, an Hplc purified non-toxic fraction derived from the bark of Quillaja Saponaria Molina. Optionally this may be admixed with 3 De-O-acylated monophosphoryl lipid A (3D-MPL), optionally together with an carrier.\nThe method of production of QS21 is disclosed in U.S. Pat. No. 5,057,540.\nNon-reactogenic adjuvant formulations containing QS21 have been described previously (WO 96/33739). Such formulations comprising QS21 and cholesterol have been shown to be successful TH1 stimulating adjuvants when formulated together with an antigen.\nFurther adjuvants which are preferential stimulators of TH1 cell response include immunomodulatory oligonucleotides, for example unmethylated CpG sequences as disclosed in WO 96/02555.\nCombinations of different TH1 stimulating adjuvants, such as those mentioned hereinabove, are also contemplated as providing an adjuvant which is a preferential stimulator of TH1 cell response. For example, QS21 can be formulated together with 3D-MPL. The ratio of QS21: 3D-MPL will typically be in the order of 1:10 to 10:1; preferably 1:5 to 5:1 and often substantially 1:1. The preferred range for optimal synergy is 2.5:1 to 1:1 3D-MPL: QS21.\nPreferably a carrier is also present in the vaccine composition according to the invention. The carrier may be an oil in water emulsion, or an aluminium salt, such as aluminium phosphate or aluminium hydroxide.\nA preferred oil-in-water emulsion comprises a metabolisible oil, such as squalene, alpha tocopherol and Tween 80. In a particularly preferred aspect the antigens in the vaccine composition according to the invention are combined with QS21 and 3D-MPL in such an emulsion. Additionally the oil in water emulsion may contain span 85 and/or lecithin and/or tricaprylin.\nTypically for human administration QS21 and 3D-MPL will be present in a vaccine in the range of 1 xcexcg-200 xcexcg, such as 10-100 xcexcg, preferably 10 xcexcg-50 xcexcg per dose. Typically the oil in water will comprise from 2 to 10% squalene, from 2 to 10% alpha tocopherol and from 0.3 to 3% tween 80. Preferably the ratio of squalene: alpha tocopherol is equal to or less than 1 as this provides a more stable emulsion. Span 85 may also be present at a level of 1%. In some cases it may be advantageous that the vaccines of the present invention will further contain a stabiliser.\nNon-toxic oil in water emulsions preferably contain a non-toxic oil, e.g. squalane or squalene, an emulsifier, e.g. Tween 80, in an aqueous carrier. The aqueous carrier may be, for example, phosphate buffered saline.\nA particularly potent adjuvant formulation involving QS21, 3D-MPL and tocopherol in an oil in water emulsion is described in WO 95/17210.\nThe present invention also provides a polyvalent vaccine composition comprising a vaccine formulation of the invention in combination with other antigens, in particular antigens useful for treating cancers, autoimmune diseases and related conditions. Such a polyvalent vaccine composition may include a TH-1 inducing adjuvant as hereinbefore described.\nWhile the invention has been described with reference to certain BASB023 polypeptides and polynucleotides, it is to be understood that this covers fragments of the naturally occurring polypeptides and polynucleotides, and similar polypeptides and polynucleotides with additions, deletions or substitutions which do not substantially affect the immunogenic properties of the recombinant polypeptides or polynucleotides.\nCompositions, Kits and Administration\nIn a further aspect of the invention there are provided compositions comprising a BASB023 polynucleotide and/or a BASB023 polypeptide for administration to a cell or to a multicellular organism.\nThe invention also relates to compositions comprising a polynucleotide and/or a polypeptides discussed herein or their agonists or antagonists. The polypeptides and polynucleotides of the invention may be employed in combination with a non-sterilc or sterile carrier or carriers for use with cells, tissues or organisms, such as a pharmaceutical carrier suitable for administration to an individual. Such compositions comprise, for instance, a media additive or a therapeutically effective amount of a polypeptide and/or polynucleotide of the invention and a pharmaceutically acceptable carrier or excipient. Such carriers may include, but are not limited to, saline, buffered saline, dextrose, water, glycerol, ethanol and combinations thereof. The formulation should suit the mode of administration. The invention further relates to diagnostic and pharmaceutical packs and kits comprising one or more containers filled with one or more of the ingredients of the aforementioned compositions of the invention.\nPolypeptides, polynucleotides and other compounds of the invention may be employed alone or in conjunction with other compounds, such as therapeutic compounds.\nThe pharmaceutical compositions may be administered in any effective, convenient manner including, for instance, administration by topical, oral, anal, vaginal, intravenous, intraperitoneal, intramuscular, subcutaneous, intranasal or intradermal routes among others.\nIn therapy or as a prophylactic, the active agent may be administered to an individual as an injectable composition, for example as a sterile aqueous dispersion, preferably isotonic.\nIn a further aspect, the present invention provides for pharmaceutical compositions comprising a therapeutically effective amount of a polypeptide and/or polynucleotide, such as the soluble form of a polypeptide and/or polynucleotide of the present invention, agonist or antagonist peptide or small molecule compound, in combination with a pharmaceutically acceptable carrier or excipient. Such carriers include, but are not limited to, saline, buffered saline, dextrose, water, glycerol, ethanol, and combinations thereof. The invention further relates to pharmaceutical packs and kits comprising one or more containers filled with one or more of the ingredients of the aforementioned compositions of the invention. Polypeptides, polynucleotides and other compounds of the present invention may be employed alone or in conjunction with other compounds, such as therapeutic compounds.\nThe composition will be adapted to the route of administration, for instance by a systemic or an oral route. Preferred forms of systemic administration include injection, typically by intravenous injection. Other injection routes, such as subcutaneous, intramuscular, or intraperitoneali can be used. Alternative means for systemic administration include transmucosal and transdermal administration using penetrants such as bile salts or fusidic acids or other detergents. In addition, if a polypeptide or other compounds of the present invention can be formulated in an enteric or an encapsulated formulation, oral administration may also be possible. Administration of these compounds may also be topical and/or localized, in the form of salves, pastes, gels, solutions, powders and the like.\nFor administration to mammals, and particularly humans, it is expected that the daily dosage level of the active agent will be from 0.01 mg/kg to 10 mg/kg, typically around 1 mg/kg. The physician in any event will determine the actual dosage which will be most suitable for an individual and will vary with the age, weight and response of the particular individual. The above dosages are exemplary of the average case. There can, of course, be individual instances where higher or lower dosage ranges are merited, and such are within the scope of this invention.\nThe dosage range required depends on the choice of peptide, the route of administration, the nature of the formulation, the nature of the subject\"\"s condition, and the judgment of the attending practitioner. Suitable dosages, however, are in the range of 0.1-100 xcexcg/kg of subject.\nA vaccine composition is conveniently in injectable form. Conventional adjuvants may be employed to enhance the immune response. A suitable unit dose for vaccination is 0.5-5 microgram/kg of antigen, and such dose is preferably administered 1-3 times and with an interval of 1-3 weeks. With the indicated dose range, no adverse toxicological effects will be observed with the compounds of the invention which would preclude their administration to suitable individuals.\nWide variations in the needed dosage, however, are to be expected in view of the variety of compounds available and the differing efficiencies of various routes of administration. For example, oral administration would be expected to require higher dosages than administration by intravenous injection. Variations in these dosage levels can be adjusted using standard empirical routines for optimization, as is well understood in the art.\nSequence Databases, Sequences in a Tangible Medium, and Algorithms\nPolynucleotide and polypeptide sequences form a valuable information resource with which to determine their 2- and 3-dimensional structures as well as to identify further sequences of similar homology. These approaches are most easily facilitated by storing the sequence in a computer readable medium and then using the stored data in a known macromolecular structure program or to search a sequence database using well known searching tools, such as the GCG program package.\nAlso provided by the invention are methods for the analysis of character sequences or strings, particularly genetic sequences or encoded protein sequences. Preferred methods of sequence analysis include, for example, methods of sequence homology analysis, such as identity and similarity analysis, DNA, RNA and protein structure analysis, sequence assembly, cladistic analysis, sequence motif analysis, open reading frame determination, nucleic acid base calling, codon usage analysis, nucleic acid base trimming, and sequencing chromatogram peak analysis.\nA computer based method is provided for performing homology identification. This method comprises the steps of: providing a first polynucleotide sequence comprising the sequence of a polynucleotide of the invention in a computer readable medium: and comparing said first polynucleotide sequence to at least one second polynucleotide or polypeptide sequence to identify homology.\nA computer based method is also provided for performing homology identification, said method comprising the steps of: providing a first polypeptide sequence comprising the sequence of a polypeptide of the invention in a computer readable medium; and comparing said first polypeptide sequence to at least one second polynucleotide or polypeptide sequence to identify homology.\nAll publications and references, including but not limited to patents and patent applications, cited in this specification are herein incorporated by reference in their entirety as if each individual publication or reference were specifically and individually indicated to be incorporated by reference herein as being fully set forth. Any patent application to which this application claims priority is also incorporated by reference herein in its entirety in the manner described above for publications and references."}
{"text": "This invention relates to coated abrasive grains and to a manufacturing method therefor, and particularly to coated abrasive grains obtained by an immersion method. More specifically, it relates to coated abrasive grains and a manufacturing method therefor in which a molten salt bath is employed, and the abrasive grains are coated with a layer of a metal compound such as at least one substance selected from metallic carbides, metallic borides, and metallic nitrides.\nAbrasive grains such as diamond abrasive grains and hard BN abrasive grains are referred to as ultrahard abrasive grains and are widely used as starting materials for abrasive articles or as abrasives. In particular, abrasive grains comprising natural or synthetic diamonds are bonded together with a binder in the form of a metal, a synthetic resin, or a vitreous inorganic material to form abrasive articles, such as grinding wheels. These abrasive articles have a wide range of uses including the grinding of carbide tools and hard materials such as ceramics, ferrite, and glass, and the cutting of stones and concrete. As a result of the spectacular development of ceramics in recent years, diamond wheels and hard BN wheels, which enable high-precision grinding of ceramics, have been used in large quantities.\nHowever, presently-used abrasive grains have the following problems.\nFor example, in the case of resin-bonded wheels or metal-bonded wheels which have diamond abrasive grains as their main component, in order to increase the bond strength between the abrasive grains and the resin or metal which serves as a binder, the surfaces of the abrasive grains are coated with a metal such as nickel. This metal coating is said to have the effect of preventing the deterioration of the resin due to heat which is generated along the edges of the abrasive grains during use of a grinding wheel. However, the metal coating has the problem that it causes the clogging of the wheel.\nOn the other hand, vitrified-bonded wheels which also have diamond abrasive grains as their main component use a ceramic as a binder. As these wheels are manufactured by firing at high temperatures, the bonding strength between the abrasive grains and the binder is high due to a solid phase reaction. However, during the high-temperature firing, the diamond abrasive grains undergo thermal corrosion, and as a result manufacture is difficult.\nJapanese Published Unexamined Patent Application No. 55-162499 (1980) disclosed a method of coating the surface of diamond abrasive grains with Ti(C,N,O) by chemical vapor deposition in order to prevent oxidation of the abrasive grains during manufacture of wheels, improve the adhesion to coatings, and prevent the clogging of the wheel during grinding. However, the adhesion of a coating formed by this method is said to be inadequate.\nThus, as the advantages of coated abrasive grains have been recognized, various methods have been proposed for improving the properties thereof. However, as all of these methods employ CVD or other vapor phase deposition processes, they are unsuitable for mass production, and they require large and expensive equipment in order to be carried out.\nThe following U.S. patents are issued regarding coated diamond abrasive grains of the prior art: U.S. Pat. Nos. 3,465,416; 3,520,667; 3,617,346; 3,650,714; 3,826,630; 3,924,031; 3,929,432; 4,063,907; 4,220,455; 4,417,906; and 4,606,738.\nAccordingly, in order to take advantage of the above-described excellent properties of coated abrasive grains, there is a need for a method which can easily and greatly improve the adhesion to coatings of coated abrasive grains and which at the same time is inexpensive."}
{"text": "1. Field of the Invention\nThis invention relates to a radiation image read-out apparatus for performing a radiation image read-out operation by use of a photomultiplier.\n2. Description of the Related Art\nIt has heretofore been proposed to use stimulable phosphors in radiation image recording and reproducing systems. Specifically, a radiation image of an object, such as a human body, is recorded on a sheet provided with a layer of the stimulable phosphor (hereinafter referred to as a stimulable phosphor sheet). The stimulable phosphor sheet, on which the radiation image has been stored, is then exposed to stimulating rays, such as a light, which cause the stimulable phosphor sheet to emit light in proportion to the amount of energy stored on the stimulable phosphor sheet during exposure of the stimulable phosphor sheet to the radiation. The light emitted by the stimulable phosphor sheet, upon stimulation thereof, is photoelectrically detected and converted into an electric image signal by use of image read-out means, such as a photomultiplier (hereinbelow referred to as the PMT). The image signal is then processed and used for the reproduction of the radiation image of the object as a visible image on a recording material.\nAs for the PMT utilized in the radiation image recording and reproducing systems described above, the sensitivity (i.e., the level of an output signal with respect to a light intensity of incident light) of the PMT is determined by the level of an applied high voltage. However, even though PMT's are of an identical type, the characteristics vary for respective PMT's. The level of the high voltage required to obtain a desired sensitivity varies for the respective PMT's. Therefore, at the time of delivery of each of the radiation image recording and reproducing systems to a user, the level of the high voltage has heretofore been set appropriately for each of the radiation image recording and reproducing systems such that a desired sensitivity of the PMT is capable of being obtained. (The setting of the level of the high voltage is described in, for example, Japanese Unexamined Patent Publication Nos. 3(1991)-132644 and 2001-100340.)\nHowever, recently, it has been found that the sensitivities of the PMT's become low due to the use of the PMT's. In particular, as for the PMT's utilized in the radiation image recording and reproducing systems described above, it has been found that, since a large anode current flows through each of the PMT's, the sensitivities of the PMT's become markedly low. FIG. 4 is a graph showing a relationship between a cumulative dose of radiation having been delivered to a stimulable phosphor sheet, i.e. a cumulative addition value of anode currents flowing through a PMT, and a sensitivity of the PMT. It has been experimentally found that, as illustrated in FIG. 4, in cases where the cumulative dose of the radiation becomes large, the sensitivity of the PMT becomes as low as approximately one tenth of the original sensitivity.\nTherefore, in cases where the voltage applied across the PMT is set appropriately at the time of the delivery of the PMT to the user, the sensitivity of the PMT becomes low due to the use of the PMT. Accordingly, particularly, appropriate detection is not capable of being made with respect to the light having a low light intensity, which light is emitted from a site of the stimulable phosphor sheet having been exposed to a small dose of the radiation. Also, in cases where a visible image is reproduced from an image signal having been obtained with the PMT whose sensitivity has become low as described above, it is not always possible to discriminate only through the seeing of the reproduced visible image whether the image quality of the reproduced visible image is due to the lowering of the sensitivity of the PMT or is due to the original quality of the radiation image having been stored on the stimulable phosphor sheet. In such cases, there is the risk that an appropriate image diagnosis will not capable of being made.\nFurther, it is considered that the lowering of the sensitivity of the PMT depends upon the cumulative value of anode currents, which have flowed through the PMT. Therefore, the extent of the lowering of the sensitivity of the PMT is not capable of being detected accurately in accordance with only the period, for which the PMT has been used, or the number of times of the radiation image read-out operations, which have been performed by use of the PMT. (The aforesaid problems with regard to the detection of the extent of the lowering of the sensitivity of the PMT are described in, for example, Japanese Patent Publication No. 7(1995)-031373.)"}
{"text": "1. Field of the Invention\nThe present invention relates to network interface devices for interconnecting host processors with a communication network, and more particularly to the processing of specific types of packets using high speed, and/or embedded hardware at the network interface.\n2. Description of Related Art\nNetwork interface cards (NIC) are often be used for network connection between a host bus system and a network medium. NICs have integrated circuits with buffer memory for storing packets in transit. For outgoing traffic, the packets being downloaded from host memory via the host bus system go through a transmit packet data buffer and medium access control MAC unit to be transmitted onto the wire, optical fiber, or antenna (or otherwise transmitted on a network medium). For incoming traffic, the packet being received from the network medium can go through the MAC unit and a receive packet data buffer to be uploaded into host memory.\nThere are times that the packets can be processed on the NIC, rather than after delivery to other processing resources. For example, in embedded firewall applications, a certain field, or fields, of the packet can be inspected to determine whether further transfer of the packet should be allowed or denied.\nManagement of computer networks is accomplished in many systems by a central network management station which has access to end stations in the network for management functions. Several specialized control packets have been developed, which are transmitted to the end stations in support of these management functions. Some of these control packets are suitable for processing, or are designed to be processed, at the network interface, rather than after delivery to the host system on which the network interface is attached.\nIn one prior art system, network interface devices are configured to capture packets while the host system is not active, including “wake up” packets using resources on the interface card. See, NIC Device-Class Power Management Specification, Version 1.0a, Nov. 21, 1997; Microsoft Corporation. (See, http://www.microsoft.com/hwdev/specs/PMref/PMnetwork.htm). The NIC Device-Class Power Management Specification handles the situation in which a host processor running Windows or another operating system OS wants to go to sleep, yet allow others to access any shared directories or devices it might have offered to the network. So the host OS passes the adapter a set of filters (where a filter in this example is a bit mask to specify which bytes in the packet are interesting, with a byte string for comparing with the interesting bytes) which the adapter should use. If a packet comes in that matches the filters, then the adapter wakes up, and signals power management resources in the host system.\nAlso, inspection of a field or fields in a packet may be required in support of logical processing of the packet for priority determinations, or other network management and system management functions.\nIn other examples, the packets processed on the NIC include a remote control packet allowing a management console to remotely reboot the targeted computer. Such a packet would simply be discarded by the interface processor and an action performed to reboot the computer, such as by sending a command to the host using a management interface like the SMBus (See, Smart Battery System Specifications—System Management Bus Specification, Rev. 1.0, (1995) Benchmarq Microelectronics, Inc., et al.).\nIn yet another example of NIC based packet processing, NIC might execute processes for which it needs to know the local internet protocol IP address of its host. This can be complicated if Dynamic Host Configuration Protocol DHCP, or another protocol for assigning dynamic IP addresses to devices on a network, is in use by which the IP address might change over time. By trapping the DHCP packets and examining them before passing them to the host, the interface card can track the changes in the IP address as they happen, and do it without adding any extra instructions to the critical code paths on the host which might increase CPU utilization or reduce performance.\nIn the prior art, to process data in a packet, the processing resources on the NIC have been required to read several fields in the packet in order to identify fields relevant to the processing being executed. For example, TCP/IP packets have certain fields that have a pre-specified location in the packet based on the layer two protocol. However, the IP header and TCP header have variable lengths. So in order to find the start of the TCP header, the NIC must first find the start of the IP header, and the length of the IP header to calculate the variable location of the start of the TCP header. Accordingly, packet processing requires multiple read and process cycles, which make it more difficult to keep up with the packet throughput in the NIC. Therefore, there is a chance that the buffer memory in the NIC will overflow, while processing resources on the NIC process a particular packet. If the memory overflows, packets are dropped and network performance suffers.\nWith 100 Megabit, Gigabit and higher wire speeds, it becomes even more difficult for the processor to inspect the packets, especially for small size packets. See, International Patent Publication No. WO01/37489 A1, entitled RECEIVE FILTERING FOR COMMUNICATIONS INTERFACE, invented by Connery, et al., published on 25 May 2001, (priority based on U.S. patent application Ser. No. 09/434,253) which describes techniques for improving performance of processor-based inspection of data packets in transit through a NIC, by reducing the number of packets that must be processing in the packet flow.\nIt is desirable to provide a network interface capable of processing packets, without incurring the increased costs associated with higher speed, powerful on chip, or on-board, processors. Also, it is desirable to extend the type of processes that may be executed by resources at the network interface."}
{"text": "The present invention relates to illumination devices generally, and more particularly to portable illumination devices for removable attachment to various objects. This invention also relates to methods for removably attaching portable illumination devices to various objects.\nA variety of illumination devices have been implemented to illuminate areas which receive no or insufficient light for human utilization. Most commonly, artificial means have been used to illuminate such areas. Some of the devices that have been developed to provide such portable illumination include electric or oil lanterns, self-supporting lights, and individual lighting units powered by batteries or electrical cords.\nA particular application for such illumination means is in illuminating desired areas for inspection and/or manipulation, where traditional lighting systems are not available. Relatively compact illumination devices, therefore, have become desirable for such purposes.\nFlashlights, in particular, have become very popular in such situations for their relative portability and low cost. One of the most common applications for such flashlights is in typical household maintenance and construction, though a wide variety of other flashlight applications are used on a daily basis. Often times, the flashlight operator utilizes the flashlight to illuminate an area in which tools or other implements are being utilized. To effectively manipulate such tools, the operator has traditionally needed to position the flashlight such that the flashlight illuminates the work area without requiring the user to continuously hold the flashlight. However, most portable flashlights project a relatively narrow illumination beam, thereby forcing the user to repeatedly re-direct and re-position the flashlight to illuminate desired areas. Thus, a need arose for a means of directing an illumination beam from a flashlight without requiring continuous manipulation from the user\"\"s hands.\nA number of mechanisms have been introduced which are designed to grasp conventional flashlights, and to be placed in various positions to more effectively illuminate desired areas. Examples of such devices are described in U.S. Pat. Nos. 4,214,688, 2,354,853, and 1,540,372. Such devices provide means for gripping conventional flashlights such that user manipulation is not required, thereby freeing the users hands to perform other tasks. However, such devices typically have limited uses, and may not be universally adapted to function in various applications.\nPortable illumination mechanisms proposed to date have not addressed the issue of having a lightweight device that may be removably attached to various objects such that the object may still be used for their intended purpose while having such an illumination device attached thereto. Furthermore, no such device has been proposed which incorporates an attachment means in the body of the portable illumination device itself, so that a single illumination device may be easily transferred to and from several receptacle objects.\nAccordingly, it is a principle object of the present invention to provide a portable illumination device having attachment means integrated therein for removable attachment of the illumination device to various objects.\nA further object of the present invention is to provide a portable illumination device that may be easily secured to and removed from small objects without interfering with the intended functionality of such objects.\nA still further object of the present invention is to provide a compact, portable illumination device that is removably attachable to various hand tools, including tape measures.\nA yet further object of the present invention is to provide a portable illumination device having universal attachment means integrated therein.\nAnother object of the present invention is to provide a portable illumination device which incorporates a plurality of distinct illumination means.\nA still further object of the present invention is to provide a portable illumination device which incorporates multi-directional illumination means, each of which illumination means has a distinct illuminatory purpose.\nA yet further object of the present invention is to provide a portable illumination device which incorporates a plurality of illumination means each emitting designated illumination wavelengths.\nA still further object of the present invention is to provide a portable illumination device incorporating one or more light emitting diodes (LED).\nBy means of the present invention, portable illumination is enhanced by providing an illumination device having attachment means integrated therein for removable attachment of the illumination device to various objects. Furthermore, such a portable illumination device may be utilized in a manner so as to avoid interfering with intended functionality of respective objects that the device is removably attached to. Such objects include everyday tools such as tape measures, saws, drills, nail guns, screw guns, as well as other common objects such as toolboxes, clothing belts, walls, and building studs. The portable illumination device of the present invention preferably utilizes one or more types of illumination elements, including light emitting diodes, incandescent lights, fluorescent lights, or other gas-charged lighting elements. Such illumination elements are preferably positioned in the illumination device and directed to illuminate specific target areas.\nOne embodiment of the portable illumination device of the present invention includes an outer frame having illumination means disposed at least partially therein, and attachment means for removably attaching the illumination device to desired objects, with the attachment means being formed integral with the outer frame. The device is preferably sized and configured to receive specific securing means which are connected to the desired objects. The attachment means may include one or more slots which are sized and configured to receive securing means being made up of protrusions sized and configured accordingly. The illumination means of the device may include one or more individual light emitting diodes (LED), wherein the light emitting diodes may illuminate in one or more distinct colors. The device may be particularly sized and configured to be removably attachable to a side surface of a tape measure such that the illumination device does not interfere with the functionality of the tape measure.\nThe present invention also contemplates a method of illuminating a target zone. The method includes providing a portable illumination device having attachment means formed integral therewith, and a securing means that is sized and configured to receive the attachment means. The method further includes connecting the securing means to a desired object, removably attaching the illumination device to the securing means via the attachment means, and actuating the illumination means. The illumination device of the present invention may also be utilized independent of the securing means, such that the illumination device is actuatable and directable in a variety of hand-held applications."}
{"text": "Passive optical networks are becoming prevalent in part because service providers want to deliver high bandwidth communication capabilities to customers. Passive optical networks are a desirable choice for delivering high speed communication data because they may not employ active electronic devices, such as amplifiers and repeaters, between a central office and a subscriber termination. The absence of active electronic devices may decrease network complexity and/or cost and may increase network reliability.\nFIG. 1 illustrates a network 100 deploying passive fiber optic lines. As shown at FIG. 1, the network 100 may include a central office 110 that connects a number of end subscribers 115 (also called end users 115 herein) within the network 100. The central office 110 may additionally connect to a larger network such as the Internet (not shown) and a public switched telephone network (PSTN). The network 100 may also include fiber distribution hubs (FDHs) 130 having one or more optical splitters (e.g., 1-to-8 splitters, 1-to-16 splitters, or 1-to-32 splitters) that generate a number of individual fibers that may lead to the premises of an end user 115. The various lines of the network 100 can be aerial or housed within underground conduits (e.g., see conduit 105).\nA portion of the network 100 that is closest to the central office 110 is generally referred to as an F1 region, where F1 is the “feeder fiber” from the central office 110. The F1 portion of the network 100 may include a distribution cable having on the order of 12 to 48 fibers; however, alternative implementations may include fewer or more fibers. A portion of the network 100 that includes at least one of the FDHs 130 and at least one of the end users 115 may be referred to as an F2 portion of the network 100. Splitters used in the typical FDH 130 may split incoming fibers of a feeder cable into, for example, 216 to 432 individual distribution fibers that may be associated with a like number of end user 115 locations.\nReferring to FIG. 1, the network 100 includes a plurality of break-out locations 125 at which branch cables 122 are separated out from main cable lines 120. The break-out locations 125 can also be referred to as tap locations, drop cable locations, splice locations or branch locations. The branch cables 122 can also be referred to as drop cables, drop lines, break-out cables or stub cables. The branch cables 122 are often connected to drop terminals 104 that include connector interfaces for facilitating coupling the fibers of the branch cables 122 to a plurality of different subscriber locations 115. The branch cables 122 are also often connected to FDHs 130.\nWithin the FDH 130, incoming optical fibers, from the central office 110, can be connected to outgoing optical fibers, leading to the end users 115, forming an optical signal connection. Each of the incoming optical fibers may be split into multiple intermediate fibers. Each of these intermediate fibers may also be connected to one of the outgoing optical fibers, forming an optical signal connection. The optical signal connections can be reconfigured within the FDH 130. For example, a particular incoming fiber may be initially connected to a first outgoing fiber but may be disconnected and instead connected to a second outgoing fiber. Certain unused incoming fibers and/or intermediate fibers may not be connected to any of the outgoing fibers. There is a need for an FDH that provides organization and storage for incoming and intermediate fibers that is adaptable to multiple configurations, the configuration being selected based on the intended application. The FDH 130 may be mounted in various applications and locations. There is a need for the mounting of the FDH 130 to be adaptable to conveniently match various applications and locations. The present disclosure satisfies these and other needs."}
{"text": "A flow cytometer system is provided. The system includes a flow cell for streaming a hydro-dynamically focused liquid sample past an interrogation zone, a light source positioned to irradiate the interrogation zone with a light beam, and beam shaping optics disposed between the light source and the flow cell. The beam shaping optics manipulates the light beam and creates a horizontally flat beam profile that irradiates the liquid sample at the interrogation zone. The system includes an intermediate angle scatter (IAS) light detector positioned to measure intermediate angle scatter emitted from the liquid sample when the liquid sample is irradiated. The system further includes a mask disposed horizontally across a portion of the IAS light detector, between the flow cell and the IAS light detector, and covering at least a central portion of the IAS light detector so as to block a horizontal diffraction pattern created by the flat beam profile irradiating the liquid sample."}
{"text": "1. Technical Field\nThe present invention relates to oil pumps suitable for supplying working oil to automatic transmissions in vehicles such as automobiles.\n2. The Related Art\nAn oil pump for an automatic transmission designed to reduce cavitation erosion is disclosed in Japanese Patent Kokai Publication No. 2003-161269. As disclosed therein, the oil pump includes a cast-iron pump body having a circular hollow formed in a face thereof; and a light-alloy pump cover closing the hollow within the pump body, thereby forming a gear compartment (pump chamber) therebetween. A drive gear is supported and driven by a drive shaft journaled in the gear compartment of the pump body, and a driven gear is disposed in the gear compartment so as to be rotatable eccentric to the drive gear and driven by the drive gear meshed with the driven gear; an arc-shaped suction port and an arc-shaped discharge port are located at the bottom of the gear compartment in a suction area and a discharge area, respectively. Working fluid spaces are defined between teeth on the outer circumferential surface of the drive gear and teeth on the inner circumferential surface of the driven gear and between side walls of the chamber bottom and the cover. Thus, the working fluid spaces are arrayed circumferentially around the pump chamber. An arc-shaped suction port is formed in the pump body and an arc-shaped discharge port is formed in the inner side face of the pump cover in the suction area and in the discharge area, respectively, for communication with the rotating working spaces.\nWith the oil pump disclosed in Japanese Patent Kokai Publication No. 2003-161269 (hereinafter simply referred to as the “prior art”), cavitation erosion is limited to an expected normal or tolerable level when the rotational speed of the drive gear is in a normal range of use (for example, up to 7,000 rpm). However, when the rotational speed of the drive gear is higher than that (for example, 7,500 rpm), the cavitation erosion of the pump cover greatly increases to an unacceptable level. This problem will now be described with reference to FIGS. 6 and 7.\nIn the prior art oil pump, a notch 5a is formed in the chamber side face of the pump body, i.e. in the “bottom” of the pump chamber (also see pump body 10 and chamber 11 in FIG. 1), and extends circumferentially from the front end of a discharge port 4a formed in the pump body to the rear end of a suction port 3a formed in the pump body in the suction area for the working spaces. In addition, a notch 5b formed in the pump cover 2, shorter than the notch 5a, extends circumferentially from the front end of a discharge port 4b formed in the cover to the rear end of a suction port 3b formed in the cover. When the drive gear 6a and driven gear 6b are rotated in the direction of the arrow during operation of the oil pump, working spaces R formed between the teeth of the gears 6a and 6b first come into communication with the discharge port 4a through the notch 5a. Since the working spaces R were in communication with the suction ports 3a and 3b until immediately before, the working spaces R are filled with low-pressure working oil entraining bubbles of a gas of volatiles from the working oil and air released from the working oil. Because the pressure of the working oil in the discharge ports 4a and 4b is significantly higher than that at the suction ports, when the working spaces R come into communication with the notch 5a, the high-pressure working oil in the discharge port 4a temporarily flows back from the pump body 1 toward the opposing inner side face of the pump cover 2 and into the working spaces R as indicated by an arrow f. Thus, the bubbles in the working spaces R collapse (become smaller), and the impact pressure occurring due to that collapsing causes cavitation erosion at the inner chamber side face of the pump cover in the vicinity where the bubbles collapse.\nWhen the rotational speed of the oil pump is less than or equal to a predetermined limit, a small number of bubbles are present in the working spaces R, the pressure of the working oil in the discharge ports 4a and 4b is not very high, and the rate of inflow into the working spaces R is also low. Therefore, the collapsing of the bubbles mainly occurs adjacent the bottom of the pump body 1, but is not relatively noticeable. Thus, cavitation erosion of the pump body 1 can be prevented by forming the pump of a material, such as cast iron, having a high resistance to cavitation erosion. Accordingly, the prior art technology is effective in preventing cavitation erosion when the rotational speed of the oil pump is less than or equal to the predetermined limit.\nHowever, when the rotational speed of the oil pump exceeds the predetermined limit, the pressure in the working spaces R is reduced, the volume of bubbles is increased, and the bubbles accumulate adjacent the inner circumference due to the increased centrifugal force. Moreover, the pressure of the working oil in the discharge ports 4a and 4b is increased, and the rate of inflow into the working spaces R is also increased. Accordingly, the position where the collapsing of the bubbles occurs is shifted to the inner side face of the pump cover 2, and more bubbles collapse. Since the pump cover 2 is composed of a material, such as aluminum, having low resistance to cavitation erosion, cavitation erosion occurs at the position indicated by symbol E1 at the inner side face of the pump cover 2, as shown in FIG. 7(b). Thus, gaps are formed between the pump gears 6a and 6b, and pump efficiency is reduced due to leaking of the working oil. It is believed that cavitation erosion occurs at the pump cover 2 by the above-described mechanism when the rotational speed of the oil pump exceeds the predetermined limit.\nTo solve the above-described problem, theoretically it would be possible to form the pump cover 2 of a metallic material having high resistance to cavitation erosion, e.g. aluminum with, for example, T6 heat treatment for increasing the surface strength or high-silicon aluminum alloy. However, such materials do not always solve the problem because of the large volume of bubbles in the working spaces R which are collapsed (crushed), and therefore, a material such as cast iron having high resistance to cavitation erosion is required. In such a case, the weight of the oil pump is disadvantageously increased since both the pump body 1 and the pump cover 2 are composed of cast iron. When such a heavy oil pump is installed in an automotive automatic transmission, the pump body or the pump cover of the oil pump cannot be integrated with the transmission housing which is composed of a light alloy, resulting in a complicated structure."}
{"text": "Despite the great amount of engineering attention given in the past to the problem of high speed feeding or sheets and other more or less flat articles, there is a continuing need for improved machines to overcome deficiencies which limit continuous, error free, high speed feeding. In particular, improvements are needed in reliably singulating sheets and other articles which vary in thickness or character, and for singulating them in a manner which does not adversely affect their surface character.\nWhen singulating paper sheets and other articles using typical prior art technology, for example--feed rollers in combination with retarder rollers or belts, double feeding or other problems such as blocking or curling can occur. One contribution to problems in even carefully configured equipment is that there is variation in the characteristics of articles being fed, even when they are of the same supposed character. Too frequently, operator intervention is required to adjust machines during operation. Similarly, extensive and careful set up is often required when changing a machine from processing one type article to another. Articles which have in the past caused feeding problems include articles having step change in thickness or varied surface texture, such as envelopes with flaps fed first, articles with body joints, window envelopes, plastic credit cards affixed to paper sheets, and heavily coated or printed articles, etc.\nAnother area where improvement is needed is in the retarders of singulators. Retarders wear if made of elastomer or plastic, or scrape unacceptable amounts of material from the sheet surfaces if made of abrasive ceramic. They tend to loose their effectiveness when debris, such as toner from photocopied pages, accretes on their operational surfaces. Thus, there is a need for conveniently and automatically presenting clean retard surfaces at singulator nips. Still another deficiency is that when sheets are taken away from a singulator at an increased speed, there is a tendency for slipping and rubbing on the surface of articles, either at the take-away rollers or the singulator feed rollers."}
{"text": "The invention relates to refrigeration systems generally used in large cooling applications. More particularly, the present invention relates to a system and method for monitoring the quantity of refrigerant within the refrigeration system.\nOne method of monitoring refrigerant includes placing a mechanical float within a receiver vessel of a refrigeration system. The mechanical float provides a visual indication of the level of refrigerant within the vessel. In this case, the level of refrigerant is only viewed during servicing operations. Alternatively, the mechanical float can include an electrical output signal fed to a tracking system. The tracking system generally includes a visual display and an alarm actuated when the level of refrigerant indicates a nearly empty receiver vessel. However, this method is difficult to employ in heat exchangers such as condensers.\nAnother method of monitoring refrigerant includes an infrared leak detector. The infrared leak detector includes a sensor placed on the outer surface of refrigeration system elements (e.g. receiver vessel, piping, valves, heat exchangers). By action of an air pump, the infrared detector can sample air surrounding the refrigeration system and detect refrigerant. The presence of refrigerant in the air can indicate the existence of a leak and thus trigger an alarm."}
{"text": "Shock absorbing devices are used in a wide variety of vehicle suspension systems for controlling motion of the vehicle and its tires with respect to the ground and for reducing transmission of transient forces from the ground to the vehicle. Shock absorbing struts are a common and necessary component in most aircraft landing gear assemblies. The shock struts used in the landing gear of aircraft generally are subject to more demanding performance requirements than most if not all ground vehicle shock absorbers. In particular, shock struts must control motion of the landing gear, and absorb and damp loads imposed on the gear during landing, taxiing and takeoff.\nA shock strut generally accomplishes these functions by compressing a fluid within a sealed chamber formed by hollow telescoping cylinders. The fluid generally includes both a gas and a liquid, such as hydraulic fluid or oil. One type of shock strut generally utilizes an “air-over-oil” arrangement wherein a trapped volume of gas (nitrogen or air) is compressed as the shock strut is axially compressed, and a volume of oil is metered through an orifice. The gas acts as an energy storage device, such as a spring, so that upon termination of a compressing force the shock strut returns to its original length. Shock struts also dissipate energy by passing the oil through the orifice so that as the shock absorber is compressed or extended, its rate of motion is limited by the damping action from the interaction of the orifice and the oil.\nOver time the gas and/or oil may leak from the telescoping cylinders and cause a change in the performance characteristics of the strut. While gas pressure can be readily monitored, it cannot be readily determined if a loss in gas pressure arose from leakage of gas alone or from leakage of both gas and oil, unless external evidence of an oil leak is noticed by maintenance personnel. If a low pressure condition is detected in the absence of external evidence of an oil leak, maintenance personnel heretofore would restore the gas pressure to a prescribed level by adding gas. This, however, eventually leads to degraded performance of the shock strut if oil had indeed escaped from the strut. Even if evidence of a oil leak is observed, maintenance personnel cannot easily determine how much oil remains or whether the remaining amount of oil meets specifications or is acceptable for operation.\nFunctionality and performance of a landing gear shock strut depends on its gas pressure and oil volume. To ensure that the landing gear functionality is within an accepted range, gas pressure and oil volume should be maintained within the design envelope. In the past, static measurement of gas pressure was the basis for shock strut servicing."}
{"text": "A touch panel, an input apparatus that allows a finger to move on a smooth panel to control a cursor, is frequently used with laptop computers, mobile phones, personal digital assistants (PDAs) and other electronic apparatus. Since a principle of the touch panel does not involve a mechanistic design that is applied to a mouse or a keyboard, it is rather simple to maintain the touch panel. As current electronic products are being developed to be more compact with better portability, the touch panel serving as an input apparatus is also becoming more popular.\nFor a touch panel, an electric signal is generated from sensing pressure on the touch panel touched by a finger or a conductivity change caused by a finger touch. Thus, a touch point sensed on the touch panel and a displacement of the touch point are influenced by a touch strength, a touched area, a hand tremor or a sensitivity of the touch panel. In addition, since a surge caused by electromagnetic interferences (EMI) may trigger the touch panel to detect an incorrect touch point, a misjudgment on the displacement of the touch point may also be concluded by the touch panel. For example, a part of certain sensors of the touch panel may be mistakenly determined as being touched since a voltage of those sensors is increased due to mobile phone electromagnetic waves, thus leading to a misjudgment concluded by the touch panel.\nTherefore, a method for determining a touch point movement and an associated apparatus capable of reducing noise that interferes with determination of a displacement of a touch point are keys for increasing efficiency of applications of a touch panel."}
{"text": "Communications and information technology equipment is commonly designed for mounting to racks and for housing within enclosures. Equipment racks and enclosures are used to contain and to arrange communications and information technology equipment, such as servers, CPUs, internetworking equipment and storage devices, in small wiring closets as well as equipment rooms and large data centers. An equipment rack can be an open configuration and can be housed within a rack enclosure. A standard rack typically includes front-mounting rails to which multiple units of equipment, such as servers and CPUs, are mounted and stacked vertically, and typically has a footprint of about 23 by 42 inches. The equipment capacity of a standard rack relates to the height of the mounting rails. The height is set at a standard increment of 1.75 inches, which is expressed as “U” units or the “U” height capacity of a rack. A typical U height or value of a rack is 42 U. A standard rack at any given time can be sparsely or densely populated with a variety of different components as well as with components from different manufacturers.\nMost rack-mounted communications and information technology equipment consumes electrical power and generates heat. Heat produced by rack-mounted equipment can have adverse effects on the performance, reliability and useful life of the equipment components. In particular, rack-mounted equipment housed within an enclosure is particularly vulnerable to heat build-up and hot spots produced within the confines of the enclosure during operation. The amount of heat generated by a rack is dependent on the amount of electrical power drawn by equipment in the rack during operation. Heat output of a rack can vary from a few watts per U unit of rack capacity to over 1 kW per U unit, depending on the number and the type of components mounted to the rack. Users of communications and information technology equipment add, remove, and rearrange rack-mounted components as their needs change and new needs develop. The amount of heat a given rack or enclosure can generate, therefore, can vary considerably from a few tens of watts up to about 10 kW.\nRack-mounted equipment typically cools itself by drawing air along a front side or air inlet side of a rack or enclosure, drawing air through its components, and subsequently exhausting air from a rear or vent side of the rack or enclosure. Air flow requirements to provide sufficient air for cooling, thus, can vary considerably as a result of the number and the type of rack-mounted components and the configurations of racks and enclosures. Generally, most configurations and designs of information technology equipment require cooling air to flow at a rate of about 180 cubic feet per minute (cfm) per kilowatt of power consumed such that a rack drawing 10 kW of electrical power would require an air flow rate of about 1,800 cfm.\nEquipment rooms and data centers are typically equipped with an air conditioning or cooling system that supplies and circulates cool air to rack-mounted equipment and enclosures. Many air conditioning or cooling systems, such as the system disclosed in U.S. Patent Application Publication No. U.S. 2001/0029163 A1, application Ser. No. 09/784,238, require that an equipment room or data center have a raised floor construction to facilitate the system's air conditioning and circulation functions. Referring to FIG. 1, the cooling system of the referenced application provides cool air by means of closed-loop air circulation and includes a raised floor 2 disposed above a base floor 5 of an equipment room. The raised floor 2 and the base floor 5 define an air passageway 6 into which an air cooling unit 14 delivers cool air. The air passageway 6 is connected to one portion of an equipment rack or enclosure 8 and is configured to channel cool air through the passageway to front portions of equipment 7 housed in the enclosure 8. Cool air flows across the equipment 7 into a plenum 8c and ascends through the plenum 8c into a plurality of ducts 24. Exhaust air vents from the ducts 24 into a return plenum 4. The return plenum 4 is connected to the air cooling unit 14 and is configured to deliver exhaust air to the air cooling unit 14 for cooling and subsequent recirculation to the equipment room.\nAlternatively, air cooling systems and methods use open floor tiles and floor grills or vents to deliver cool air from the air passageway disposed below the raised floor of an equipment room. Open floor tiles and floor grills or vents are typically located in front of equipment racks and enclosures, and along aisles between rows of racks and enclosures arranged side-by-side.\nThe cooling systems and methods that require a raised floor construction typically do not efficiently meet the cooling requirements of rack-mounted equipment. In particular, racks that include high power equipment having a thermal exhaust air output above 5 kW and up to 10 kW present a particular challenge for such systems and methods. A raised floor construction typically provides an open floor tile or a floor grill or vent having a venting area of about 12 by 12 inches and is configured to deliver from about 200 to about 500 cfm of cool air. The air flow rate from the venting area is dependent on such factors as static air pressure and the presence of other floor tiles such that, in practice, a floor tile typically delivers from about 100 to about 200 cfm of air. A rack of high power equipment drawing up to 10 kW and requiring an air flow of approximately 1,800 cfm, therefore, would need at least about 3.5 to about 5 open floor tiles, grills or vents disposed around the rack's perimeter to supply sufficient cool air to meet its cooling requirements. Such a floor configuration would be difficult to achieve in equipment rooms crowded with racks and enclosures, and impossible to implement if racks and enclosures are arranged side-by-side in rows. Air cooling systems and methods that incorporate raised floor configurations, thus, are typically only used with racks and enclosures spaced apart to provide sufficient floor area to accommodate multiple open floor tiles, grills or vents. For typical rack spacing, this places a limit on the density of equipment that can be achieved. In addition, such air cooling systems and methods must supply cold air through open floor tiles, grills or vents to meet the cooling requirements of equipment having high thermal exhaust air output.\nEquipment rooms and data centers are often reconfigured to meet new and/or different equipment needs that require individual racks and enclosures to be relocated and/or replaced. In this context, raised floor air cooling systems and methods are inflexible and can typically only be reconfigured and/or retrofitted to service rearranged, relocated and/or newly installed equipment racks at considerable cost. Raised floor configurations cannot easily and inexpensively accommodate the manner by which users typically deploy equipment racks and reconfigure equipment rooms and data centers to meet their new or changing needs.\nIn addition, cooling systems and methods that require raised floor construction lack physical flexibility and portability to operatively account for a wide variation in electrical power consumption between different racks and enclosures in an equipment room, and, in particular, between racks and enclosures located in the same row. Cooling systems and methods that rely upon raised floor air passageways and open floor tiles, grills or vents to supply cool air cannot easily and inexpensively vary or concentrate cool air to those high power racks that consume relatively large amounts of electrical power and have a high thermal air exhaust output. In addition, newly installed equipment may draw more electrical power than replaced or existing equipment to create thermal problem areas in functioning equipment rooms.\nFurther, cooling systems and methods that depend upon raised floor construction cannot physically accommodate specific areas in equipment rooms with thermal problems. For example, such systems and methods cannot overcome exhaust problems where hot and warm exhaust air is not effectively vented from an equipment room and/or returned to an air conditioning or cooling system without considerable expense to reconfigure and/or retrofit the air systems and/or the racks and enclosures. Such exhaust problems can also cause hot and warm exhaust air to recirculate to racks and enclosures, raising operating temperatures of equipment and, in particular, equipment with high power consumption and/or equipment located in areas of an equipment room with thermal problems. Similarly, raised floor cooling systems and methods cannot account for local thermal problems where hot and warm exhaust air from one rack is drawn into adjacent and/or proximate racks causing equipment overheating. Raised floor cooling systems and methods cannot easily retrofit to existing cooling systems of an equipment room or data center or to existing equipment racks.\nAir cooling systems and methods that require raised floor configurations also reduce available headroom in equipment rooms and data centers. Wires housed within air passageways below raised floors are difficult to access. In addition, such air passageways are difficult to clean. Open floor tiles and floor grills or vents also pose a safety risk to personnel. In addition, raised floor construction presents a risk of collapse during earthquake."}
{"text": "The present invention generally relates to blinds and windbreaks as used by hunters."}
{"text": "1. (Field of the Invention) 5 The present invention relates to the method of forming electrodes of piezoelectric ignition elements used in lighters or the like, capacitors, varistors, sensors, and hybrid integrated circuit (HIC) boards.\n2. (Description of the Prior Art)\nConventionally the electrode of Ceramic capacitors, piezoelectric devices and the like has been formed as shown in FIG. 2. In the figure, a cylindrical ceramic composition (a Pb(ZrTi)O.sub.3 material) 10 has its upper and lower planar surfaces coated by baking with frit-contained silver electrodes (a Ag material) 11. However, this conventional ceramic composition has problems of high manufacturing cost and high product cost due to the use of silver for the electrode. In addition, electromigration specific to silver (development of dendrite silver originating from the cathode electrode when a voltage is applied across two electrodes interleaved by a relatively high hygroscopic material) which grows gradually causes ultimately the breakdown of insulation. Moreover, ceramic devices produced by the conventional method lack in durability."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for controlling loading of active agents into liposomes. More particularly, the present invention relates to methods for modulating loading of active agents into multivesicular liposomes.\n2. Description of Related Art\nOptimal treatment with many drugs requires that the drug level be maintained at a specified level for a prolonged period of time. For example, optimal anti-cancer treatment with cell cycle-specific antimetabolites requires maintenance of a cytotoxic drug level for a prolonged period of time. Cytarabine is a highly schedule-dependent anti-cancer drug. Because this drug kills cells only when they are synthesizing DNA, prolonged exposure at therapeutic concentration of the drug is required for optimal therapeutic effect. The therapeutic effectiveness of such agents is often further complicated by the fact that the half-life after an intravenous or subcutaneous dose may be as short as a few hours. To achieve optimal therapeutic effect against cancer cells with a cell cycle phase-specific drug like cytarabine, there are two major requirements: first, the cancer cells must be exposed to a high concentration of the drug without doing irreversible significant harm to the host; and second, the tumor must be exposed to the drug for a prolonged period of time to maximize the number of cancer cells that are contacted during DNA synthesis, the susceptible portion of the cycle of cell proliferation. This kind of treatment regimen requires a high drug load in a slow release formulation.\nCertain other types of drugs are so toxic that it is important to maintain a low level of the drug over an extended period of time. For instance, amikacin is an aminoglycoside antibiotic with clinically significant activity against strains of both gram negative and gram positive bacteria. Under existing therapeutic procedures, the drug is normally administered by intravenous or intramuscular routes on a once or twice a day schedule. The most commonly used clinical dose is 15 mg/Kg/day, which is equivalent to a maximum recommended daily dose of 1 g per day. However, administration of the drug by spaced injections results in systemic exposure to the patients, and depending on the drug, increases risk of toxic side effects. Consequently, a local depot slow-release preparation for treatment of infections such as those confined to a local region of soft tissue or bone would be advantageous in increasing local tissue levels of the drug, compared with therapeutic systemic doses, while reducing or avoiding the systemic toxicity of the free drug. If the drug is highly toxic or the treatment regimen requires a low therapeutic dose, a relatively low drug load in a slow release formulation is beneficial.\nOne approach which has been used to provide controlled release compositions for drug delivery is liposome encapsulation. Among the main types of liposomes, multivesicular liposomes (Kim, et al., Biochim. Biophys. Acta; 728:339-348, 1983), are uniquely different from unilamellar liposomes (Huang, Biochemistry; 8:334-352, 1969; Kim, et al., Biochim. Biophys. Acta; 646:1-10, 1981), multilamellar liposomes (Bangham, et al., J. Mol. Bio., 13:238-252, 1965), and stable plurilamellar liposomes (U.S. Pat. No. 4,522,803). In contrast to unilamellar liposomes, multivesicular liposomes contain multiple aqueous chambers. In contrast to multilamellar liposomes, the multiple aqueous chambers of multivesicular liposomes are non-concentric.\nThe prior art also describes methods for producing multivesicular liposomes (Kim, et al., Biochim. Biophys. Acta, 728:339-348, 1983). However, the encapsulation efficiency of some small molecules, such as cytosine arabinoside, also known as cytarabine or Ara-C, proved relatively low, and the release rate of encapsulated molecules in biological fluids was faster than is therapeutically desirable. EP 0 280 503 B1 discloses a method developed for controlling the release rate of encapsulated molecules from multivesicular liposomes wherein a hydrochloride is introduced into the encapsulation process to control the rate of release in biological fluids (of the active agent. Further research, disclosed in WO 95/13796, has shown that the release rate of agents from multivesicular liposomes in human plasma can be controlled by introduction of a non-hydrochloride acid into the aqueous solution in which the agent is dissolved prior to forming the multivesicular liposome\nU.S. Pat. No. 5,077,056 discloses studies that show the rate of release of the encapsulated biological agent from liposomes into an aqueous environment can be modulated by introducing protonophores or ionophores into liposomes to create a membrane potential. In addition, a method is known (U.S. Pat. No. 5,186,941) for controlling the release rate of drugs from vesicle compositions wherein the liposomes containing a therapeutic agent encapsulated are suspended in a solution containing sufficient solute to provide an osmolarity substantially isotonic with respect to that of the solution within the vesicles, and hypertonic with respect to physiological saline. In multivesicular liposomes, it is also known (WO 96/08253) to control the rate of release of active agents by introducing an osmotic spacer into the aqueous solution in which the active agent is dissolved prior to formation of the multivesicular liposomes.\nIn addition to the biologically active agent and acids or osmotic spacers intended to control the rate of release of the biologically active agent from the liposomes, it is common practice to coencapsulate compounds that are intended to serve any of a number of helper functions. For instance, certain biologically active compounds retain activity only when kept at a particular pH. Thus acids or buffers are often necessarily encapsulated in addition to the active agent to control the pH of the drug environment. In other cases, a counterion is incorporated to enhance solubility of a biologically active agent that has a low solubility.\nThese methods for producing liposome formulations with slow release characteristics have sometimes proven incompatible with the goal of producing liposomes containing a high load of active agent with good encapsulation efficiency so that little of the expensive active agent is wasted by failure to capture it within the liposomes.\nThus the need exists for new methods for producing liposomes, for instance multivesicular liposomes (MVLs), that allow for control of drug loading, either high or low, while maintaining desirable slow release of the active agent into storage and biological fluids. Of particular interest is the development of high-load, controlled release formulations for peptides and proteins. A need also exists for new methods of achieving these goals without sacrifice of high encapsulation efficiency to avoid the waste of expensive active agents, such as drug and therapeutic proteins."}
{"text": "As more women become aware that breastfeeding is the best source of nutrition for a baby, and also offers health benefits to the nursing mother, the need is increasing for breast pump solutions that are user-friendly, quiet, discrete and versatile for use by a nursing mother in various situations. This is particularly true for the working mother, who is away from the home for eight to ten hours or more and needs to pump breast milk in order to have it available for her baby, but it is also a requirement for many other situations where the mother is away from the privacy of the home for an extended period, such as during shopping, going out to dinner or other activities.\nAlthough a variety of breast pumps are available, most are awkward and cumbersome, requiring many parts and assemblies and being difficult to transport. Hand pump varieties that are manually driven are onerous to use and can be painful to use. Some powered breast pumps require an AC power source to plug into during use. Some systems are battery driven, but draw down the battery power fairly rapidly as the motorized pump continuously operates to maintain suction during the milk extraction process. Many of the breast pumps available are clearly visible to an observer when the mother is using it, and many also expose the breast of the mother during use.\nThere is a continuing need for a small, portable, self-powered, energy efficient, wearable breast pump system that is easy to use and is discrete by not exposing the breast of the user and being invisible or nearly unnoticeable when worn."}
{"text": "The present invention is related to improved railroad car brake control valves and particularly to such improvements as relate to the adverse effects of vibration on piston valves and the like.\nRailroad car brake control valves, such as ABD, ABDW and ABDX valves, presently employ piston operated valves comprising a piston member that operates in a bushed bore and a slide valve member that is carried by the piston member for controlling the pneumatic porting between the slide valve face and the mating face of a slide valve seat formed in the bushing. The slide valve/seat interface requires a very close, lapped-fit. A special lubricant is provided to assure controlled, low friction movement of the slide valve and piston. Diaphragm operated pistons are employed, and a close fit is also provided between the piston member guides and the generally cylindrical piston bush in order to support and guide the piston member in the bushing bore.\nWith the trend toward light-weight railroad cars, such as aluminum coal hopper cars, for example, relatively severe vibration can occur when such light-weight cars are operated under less than ideal dynamic conditions. Wheel surface irregularities including out-of-roundness, worn truck components, and rail conditions are believed to contribute to such dynamic operating conditions that result in severe vibration of the brake control valve on the car.\nUnder conditions in which severe vibration is encountered, vibration can be transmitted to the piston member. Presently, the piston member employed in the industry standard ABD, ABDW, and ABDX type brake control valves is made of die cast aluminum, while the bushing in which the piston member operates is brass. Such high levels of vibration as are sufficient to induce vibration of the piston member can cause excessive wear of the piston member guide due to progressive abrasion. Such abrasive wear of the aluminum piston member has been found to cause general contamination of the piston and slide valve with a fine, oxidized aluminum powder resulting from the accumulation of particulates of aluminum formed by the abrasion of the piston guides. This residue tends to absorb the slide valve lubricant which also increases slide valve friction and causes high piston operating pressure differentials and consequent erratic and unacceptable behavior of the aforementioned control valves."}
{"text": "Conventionally, in a manufacturing process of a semiconductor element that requires a crystal film having a relatively large film thickness, like a power device such as an IGBT (Insulated Gate Bipolar Transistor), an epitaxial growth technique of performing vapor phase growth of a single-crystal thin film on a substrate such as a wafer for film formation is employed.\nIn a film formation apparatus used for the epitaxial growth technique, for example, a wafer is placed in a film formation chamber that is kept at an ordinary pressure or a reduced pressure. While this wafer is heated, a gas serving as a source for film formation (hereinafter, also simply “source gas”) is supplied to the film formation chamber. Subsequently, a thermal decomposition reaction and a hydrogen reduction reaction of the source gas occur on a surface of the wafer, thereby forming an epitaxial film on the wafer.\nTo manufacture an epitaxial wafer having a large film thickness with a high yield, it is necessary for a surface of a wafer to successively contact new source gases to increase the vapor-phase growth rate. Therefore, epitaxial growth is performed while a wafer is rotated at a high speed (see, for example, Patent Document 1).\nIn conventional film formation apparatus, a gas supply unit that supplies a source gas is provided at the top of a film formation chamber. A shower plate having a large number of source gas ejection holes formed thereon is connected to the gas supply unit. By using this shower plate, the source gas flow in the film formation chamber is uniformized and the source gas is uniformly supplied onto a wafer.\nBecause high-temperature heating at 1600° C. or more is required in SiC epitaxial growth and the like, a film formation apparatus including an auxiliary heater that heats a wafer from above has been used in addition to a heater that heats a wafer from below. This type of film formation apparatuss performing high-temperature heating have a problem that a supplied source gas contacts an inner wall surface of a film formation chamber, a film is deposited on the inner wall surface, and the film deposited on the inner wall surface is peeled off, so that particles are generated. There is another problem that, when an etching gas is supplied from a gas supply unit to clean a deposited material, an inner wall member of the film formation chamber is corroded, so that particles are generated. Furthermore, there is still another problem that it takes a long time to increase the temperature within the film formation chamber or to cool the film formation chamber.\nAn object of the present invention is to provide a film formation apparatus and a film formation method that can suppress generation of particles and can reduce a temperature-increasing time and a cooling time."}
{"text": "1. Field of the Invention\nThis invention relates to electrosurgery, and more particularly to a new and improved electrosurgical technique for achieving coagulation or a hemostatic effect, i.e. fulguration and desiccation, by conducting radio frequency (RF) electrical energy through a conductive inert gas stream to the tissue. In addition, the present invention relates to an electrosurgical fulguration arcing technique of creating an eschar and tissue effects offering a substantially improved capability for coagulation. Further still, the present invention relates to an electrosurgical non-arcing desiccation technique and equipment for applying electrical energy to tissue to achieve superior thermal desiccative effects.\n2. Description of the Prior Art\nElectrosurgery involves the application of radio frequency electrical energy to tissue. The electrical energy originates from an electrosurgical generator (ESG) and is applied by an active electrode to the tissue. The active electrode typically has a small cross-sectional or limited surface area to concentrate the electrical energy at the surgical site. An inactive return electrode or patient plate contacts the patient at a remote location from the surgical site to complete the circuit through the tissue to the ESG. The patient plate is relatively large in size to avoid destructive energy concentrations. Alternatively, a pair of active electrodes may be used in a \"bipolar\" mode in which the electrosurgical energy flows directly through the tissue between the two active electrodes, and the electrosurgical effects are confined to the tissue directly located between the two closely-spaced electrodes.\nA variety of different electrosurgical effects can be achieved, depending primarily on the characteristics of the electrical energy delivered from the ESG. Among the effects are a pure cutting effect, a combined cutting and hemostasis effect, a fulguration effect and a desiccation effect. Desiccation and fulguration are usually described collectively as coagulation. Many conventional ESG's offer the capability to selectively change the energy delivery characteristics and thus change the electrosurgical effects created.\nSatisfactory fulguration effects have been particularly difficult to obtain. Some surgeons have preferred to use older spark gap generators known as \"Bovie\" devices for fulguration, but use other more modern ESG's for cutting or cutting with hemostasis. Indeed, spark gap ESG's have been the standard against which modern solid state ESG's have been measured for achievement of satisfactory fulguration effects. One modern ESG which achieves substantially improved fulguration effects, compared to both spark gap and previous solid state ESG's is described in U.S. Pat. No. 4,429,694, assigned to the assignee hereof. Despite the improvements available in fulguration, certain disadvantages remain for which there have been no satisfactory alternatives.\nConventional fulguration is characterized by electrical arcing through the air from various locations on the metal surface of the active electrode, with the arcs contacting the tissue in somewhat of a random non-predictable manner. In many cases, arcs leave the active electrode in an initial trajectory traveling away from the tissue before actually curving around and striking the tissue surface. The result is an uneven, randomly concentrated or distributed delivery of arcing energy. An uneven eschar of variable characteristics is created on the surface of the tissue.\nThe random delivery of the arc energy creates holes which are significantly disparate in diameter (or cross-sectional size) and in depth. Larger, deeper holes are formed by repeated arcs contacting the tissue at approximately the same location. Smaller arc holes are also present in the tissue but they are unevenly distributed about the larger arc holes. The smaller arc holes are created by single individual arcs, or the less repetitious arcing to the tissue at the same location. The smaller arc holes are relatively small in diameter or cross-section and relatively shallow in depth, compared to the larger arc holes. Significant variations in cross-sectional size and depth between the large and small arc holes occur. Significant variations exist in the spacing and in the amounts of tissue between the large and small arc holes, causing the substantial variations in the surface distribution of the holes.\nThermal necrosis occurs in the tissue between the arc holes. The degree of thermal necrosis varies between total carbonization between the more closely spaced larger holes, to necrosis without charring or carbonization between the more widely separated smaller arc holes.\nThe eschar created has two distinct layers above the unaffected viable tissue. An arc hole reticulum of the tissue subjected to necrosis is created by the pattern of arc holes. The arc hole reticulum extends to greater depths in the areas of the deeper arc holes, and to substantially shallower depths in the areas of the shallower arc holes. Due to the random distribution and depth of the arc holes, the arc hole reticulum is relatively uneven in depth. Significant variations in the depth of the arc hole reticulum layer are typical. A layer of thermally desiccated tissue is located below the arc hole reticulum layer. Tissue necrosis in the thermally desiccated layer occurs as a result of the current heating effects of the electrical energy dissipating from the arcs. The desiccation layer is also uneven in depth and location due to the nonuniform application of the arcing energy over the arc hole reticulum layer. Significant variations in the depths of the desiccation layer are also typical.\nOver a given area of tissue, certain locations are only moderately affected by the arcing energy. A thin arc hole reticulum and a thin desiccation layer result. Other areas have a relatively thick eschar formed therein. Very thick carbonized eschars tend to be fragile and are prone to crack when flexed, usually resulting in renewed bleeding from the unaffected tissue below the desiccation layer. Thin eschars are more flexibile and therefore more desirable, but it has been difficult to obtain sufficient coagulation effects from thin eschars.\nCauses of the uneven eschar created by prior fulguration techniques are not known with certainty, but numerous factors are theorized to play a role. One of the more significant contributory factors is probably changes in impedance in the arc pathway between the active electrode and the tissue. Impedance changes may result from variations in the distance which the arcs travel through the air, due to the changes in ionization potential between the active electrode and the tissue. It is virtually impossible for the surgeon to maintain the active electrode at a consistent distance from the tissue, particularly if the tissue is moving due to pulsation, or due to puckering and swelling as a result of applying the electrical energy. The arcing from random locations on the active electrode also creates different arc length pathways and hence impedances. The combined impedance of the tissue and the eschar changes with the application of electrical energy. The volatilization of the cells and vaporization of the moisture in the cells changes the relative impedance in a localized spot-to-spot manner on the surface of the tissue. The formation of the charred material also influences the arc pathways, presenting an opportunity for subsequent arcs to return to the tissue at the same location and thereby enlarge the pre-existing arc hole and create even further charring.\nAnother problem with conventional electrosurgery is that it is very difficult if not impossible to achieve effective fulguration on spongy or vascular tissue such as the liver or the spleen, or on other tissues from which there is a tendency for blood to continually ooze over the surface from the highly developed vascular network within the tissue. Often, only the surface of the oozing blood is coagulated, with no penetration to the surface of the tissue below the layer of blood. A superficial coagulum results on the surface of the blood, but this coagulum quickly sloughs away resulting in only temporary hemostasis. Of course, once the temporary coagulum sloughs away, bleeding continues. Even if a coagulation effect on the tissue surface can be established, it is easily destroyed or perforated by the arcs returning to the same locations causing the longer, deeper arc holes. The deeper arc holes perforate the eschar and extend into the viable tissue below the eschar to provide a pathway for continued bleeding. The heat created by the arcs causes boiling of moisture below the eschar, and the pressure of resulting vapor can also rupture the eschar to reinitiate bleeding.\nApart from the tissue disadvantages of conventional electrosurgical fulguration, certain other practical problems exist. Arcing from the active electrode rapidly increases the temperature of the active electrode. Electrode heating is responsible for a number of problems. If the heated active electrode contacts the tissue, as it inevitably will, or if the active electrode is immersed in fluid such as blood, proteins from the tissue or the blood are denatured and stick to the active surface of the electrode. The buildup of charred material on the electrode eventually creates a sufficiently high impedance so that adequate power can no longer be delivered. The surgeon must continually clean the electrode by wiping or scraping the charred material, which disrupts, distracts, and prolongs the surgical operation. Freshly created eschars can be detached in an effort to free a sticking electrode from the tissue surface. The random accumulation of charred material on the active electrode creates more random delivery of the arcing energy, even further increasing the random delivery pattern. Because of the variable nature of the impedance of the charred material, consistent power application is difficult or impossible. The accumulation of the charred material can obscure the surgeons view of the surgical site. The temperature of the active electrode may reach sufficiently elevated levels to transfer molten metal from the electrode to the patient, creating questionable effects. Because the electrode contacts the tissue, there is a potential for cross-contamination between viable tissues and diseased tissues. Although the clinical problems associated with cross contamination are not fully understood at the present time, the advantages of eliminating the possibility are evident. A significant smoke plume also results from the burning tissue because of the air environment in which the electrosurgery occurs. Not only does the plume produce a noxious odor, but there may be some evidence that particulates in the smoke plume from burning tissue may contain hazardous chemicals, virus, bacteria, neoplastic cells and other hazards. Of course, the oxygen environment in which the electrosurgery is conventionally conducted exhibits a potential for igniting paper drapes, surgical sponges and the like.\nSome of the typical problems associated with creating and applying the arcs in conventional electrosurgery can be improved by optimizing the operating and other characteristics of the electrosurgical generator. U.S. Pat. No. 4,429,694 discloses an improved ESG which reduces some of the described disadvantages during fulguration. However, many of the disadvantages cannot be avoided and many of the characteristics cannot be improved by conventional electrosurgical techniques and equipment, due to the limitations previously inherent in electrosurgery.\nThe conventional technique of obtaining thermal desiccation by use of a conventional ESG is to apply electrical energy from a flat surface of the active electrode placed in contact with the tissue. An electrical resistance heating effect is created by the current flowing into the tissue from the active electrode. Because the active electrode contacts the tissue surface over a relatively large area, no arcing is intended to occur. To spread the thermal desiccation effect over a substantially large area, the active electrode is moved from location to location. It is very difficult to apply a level of energy which will obtain thermal desiccation but which will not cause the tissue to stick on the flat surface of the active electrode or arcing from the active electrode to non-contacted surface areas. The thermal desiccation effects are unevenly distributed because the active electrode is moved from spot to spot. Overlapping the spots of energy application can enhance the probability for tissue sticking and exaggerate the variable depth effects. Of course, moving the active electrode from spot to spot is very time consuming in an operation where time is very important or critical.\nThe prior art desiccation technique can only be applied to create surface desiccation effects. Furthermore, the inability to accurately control the amount of power, tissue sticking effects, and the like have prevented the prior use of electrosurgery on very thin fragile tissue such as the mesentary, and in other surgical techniques.\nIt is against this abbreviated background of previously existing disadvantages and problems in electrosurgery that the advantages and improvements of the present invention can be better appreciated."}
{"text": "1. Field of the Invention\nThe present invention relates to a data reproduction protection method, and a data reproduction apparatus for implementing such a protection method, whereby reproduction of a signal represented by digital data such as a recorded digital video signal can be selectively restricted.\n2. Description of the Prior Art\nIn the following, the term xe2x80x9cdata mediumxe2x80x9d is to be understood in a very general sense, as applying for example to broadcasting systems which transmit data such as video and/or audio data, in which case the received data may be the object of reproduction protection, and as applying also to any type of recording medium such as recording disks or tapes, etc., in which case playback data derived from the recording medium may be the object of reproduction protection. The reproduction protection may serve to selectively restrict viewing, hearing or copying of the data.\nIn the prior art, various types of reproduction protection method have been applied in fields such as CATV (cable television) and satellite television broadcasting. One method is to execute scrambling processing of transmitted video and audio data, and to insert a copyright code into the data, for thereby dividing the data into portions which can be freely reproduced and portions for which a fee must be paid in order to reproduce the data. When a program for which payment of a fee is necessary is received by a receiving apparatus, the program can be unscrambled and reproduced only if specified payment conditions are satisfied.\nIn the case of recorded media, one method of reproduction protection which is applicable to the DAT (digital audio tape recorder) recording system is the SCMS (serial copy management system). With that method, the playback DAT signal from a DAT playback apparatus has a main ID (identification) number which includes a copy inhibit code, whereby a single [copy enablexe2x80x94copy inhibit] sequence is ensured, so that a user can only make a single copy of a pre-recorded digital audio tape.\nHowever with such prior art methods of reproduction protection there are only two control possibilities, i.e. reproduction is made either possible or impossible. It has not been possible hitherto to provide a gradually varying degree of restriction of reproduction of a signal conveyed by a data medium. Thus, such a reproduction protection method can only be used for a single purpose, e.g. for management of payment fees, or for copyright protection. Moreover with such a prior art reproduction protection method, since the data which are to be protected exist only in a transmitting medium or recording medium prior to being reproduced, it has not been possible to provide a varying degree of limitation of reproduction capability in accordance with some condition of the reproduction apparatus. Thus in some cases, the degree of protection may be excessively severe, or excessively lax, so that it is difficult to achieve an effective degree of protection. For example, certain types of scenes recorded on a video tape may be permitted to be viewed in a certain country, such as the U.S.A., but may not be permissible in other countries. It would thus be advantageous to ensure that when that video tape is played on a reproduction apparatus which is sold to the public in such other countries, reproduction protection is automatically applied such that the aforementioned scenes will not be reproduced, or will not be clearly reproduced. However in the prior art, such a feature has not been possible.\nIt is an objective of the present invention to overcome the problems of the prior art set out above, by providing a reproduction protection method and apparatus whereby information specifying a degree of restriction of reproduction of an original signal is conveyed (e.g. by a recording medium or signal transmission medium) together with data expressing the original signal, whereby information specifying a degree of restriction of reproduction of the original signal are generated by a reproduction apparatus which operates on the conveyed data, and whereby information specifying a degree of restriction which is actually applied to reproduction of the original signal is derived based on a combination of the restriction information conveyed by the data medium and the restriction information generated by the reproduction apparatus.\nMore specifically, the invention provides a reproduction protection method comprising: attaching medium protection data to main data which are conveyed by a data medium, said main data representing an original signal;\nsupplying the main data and medium protection data, via the data medium, to a reproduction apparatus;\ngenerating apparatus protection data by the reproduction apparatus;\ndetermining a protection level by combining the medium protection data and the apparatus protection data; and\ncontrolling the reproduction apparatus to utilize the main data to reproduce the original signal in accordance with the protection level.\nIt is a further objective of the invention to overcome the above problems by providing a reproduction apparatus providing reproduction protection, for operating on main data representing an original signal and medium protection data expressing a medium protection level, said main data and medium protection being conveyed by a data medium, the apparatus comprising:\nmeans for detecting said medium protection data to obtain a medium protection signal expressing said medium protection level;\nmeans for generating an apparatus protection signal expressing an apparatus protection level which has been assigned to said reproduction apparatus;\nmeans responsive to said medium protection signal and apparatus protection signal for determining a final protection level in accordance with a combination of said medium protection level and apparatus protection level;\nmeans for executing reproduction of said original signal by utilizing said main data, including means for selectively restricting said reproduction in accordance with said final protection level.\nWith such a method and apparatus for reproduction protection, the protection level can be determined in accordance with the medium protection data, and hence can be determined in accordance with the wishes of the manufacturer of the data medium, or of the copyright owner of the main data. In addition, the protection level which is actually applied (i.e. the final protection level) is also determined in accordance with the apparatus protection data, which can be specified by the manufacturer or the seller of the reproduction apparatus. As a result, when the main data are to be reproduced (for example, during playback of a recording disk or tape), a graduated degree of limitation of reproduction is implemented, with that degree of limitation being determined by the final protection level, i.e. being determined in accordance with a combination of the requirements of the data medium manufacturer or the copyright owner of the main data and the requirements of the manufacturer or seller of the reproduction apparatus. In that way, considerable flexibility can be ensured in selectively restricting reproduction of signals which are conveyed for example by recording disks or tapes or by broadcasting systems.\nIn particular, the invention enables such reproduction restriction to be applied to specific frames or sequences of frames of a video signal, or to specific regions within each of a sequence of frames."}
{"text": "It is known to provide an inflatable vehicle occupant protection device, such as an air bag, for helping to protect an occupant of a vehicle. A driver frontal air bag is inflatable between an occupant of a front seat of the vehicle and the steering wheel of the vehicle. When inflated, the air bag helps protect an occupant from impacts with parts of the vehicle, such as the steering wheel of the vehicle. During inflation, it may be desirable to prevent immediate contact between the air bag and inflation fluid exiting an inflator."}
{"text": "Two of the most common types of osteoporosis are postmenopausal and senile osteoporosis. Osteoporosis affects men as well as women, and, taken with other abnormalities of bone, presents an ever-increasing health risk for an aging population. The most common type of osteoporosis is that associated with menopause. Most women lose between 20-60% of the bone mass in the trabecular compartment of the bone within 3-6 years after the cessation of menses. This rapid loss is generally associated with an increase of bone resorption and formation. However, the resorptive cycle is more dominant and the result is a net loss of bone mass. Osteoporosis is a common and serious disease among postmenopausal women. There are an estimated 25 million women in the United States alone who are afflicted with this disease. The results of osteoporosis are both personally harmful, and also account for a large economic loss due to its chronicity and the need for extensive and long-term support (hospitalization and nursing home care) from the disease sequelae. This is especially true in more elderly patients. Additionally, while osteoporosis is generally not thought of as a life-threatening condition, a 20-30% mortality rate is related to hip fractures in elderly women. A large percentage of this mortality rate can be directly associated with postmenopausal osteoporosis.\nThe most vulnerable tissue in the bone to the effects of postmenopausal osteoporosis is the trabecular bone. This tissue is often referred to as spongy bone and is particularly concentrated near the ends of the bone near the joints and in the vertebrae of the spine. The trabecular tissue is characterized by small structures which inter-connect with each other as well as the more solid and dense cortical tissue which makes up the outer surface and central shaft of the bone. This cris-cross network of trabeculae gives lateral support to the outer cortical structure and is critical to the biomechanical strength of the overall structure. In postmenopausal osteoporosis, it is primarily the net resorption and loss of the trabeculae which lead to the failure and fracture of the bone. In light of the loss of the trabeculae in postmenopausal women, it is not surprising that the most common fractures are those associated with bones which are highly dependent on trabecular support, e.g., the vertebrae, the neck of the femur, and the forearm. Indeed, hip fracture, Colle's fractures, and vertebral crush fractures are indicative of postmenopausal osteoporosis.\nOne of the earliest generally accepted methods for treatment of postmenopausal osteoporosis was estrogen replacement therapy. Although this therapy frequently is successful, patient compliance is low, primarily due to the undesirable side-effects of chronic estrogen treatment. Frequently cited side-effects of estrogen replacement therapy include reinitiation of menses, bloating, depression, and fear of breast or uterine cancer. In order to limit the known threat of uterine cancer in those women who have not undergone a hysterectomy, a protocol of estrogen and progestin cyclic therapy is often employed. This protocol is similar to that which is used in birth control regimens, and often is not tolerated by women because of the side-effects characteristic of progestin. More recently, certain antiestrogens, originally developed for the treatment of breast cancer, have been shown in experimental models of postmenopausal osteoporosis to be efficacious. Among these agents is raloxifene (See, U.S. Pat. No. 5,393,763, and Black et al, J. Clin. Invest., 93:63-69 (1994)). In addition, tamoxifene, a widely used clinical agent for the treatment of breast cancer, has been shown to increase bone mineral density in post menopausal women suffering from breast cancer (Love et al, N. Engl. J. Med., 326:852-856 (1992)).\nAnother therapy for the treatment of postmenopausal osteoporosis is the use of calcitonin. Calcitonin is a naturally occurring peptide which inhibits bone resorption and has been approved for this use in many countries (Overgaard et al, Br. Med. J, 305:556-561 (1992)). The use of calcitonin has been somewhat limited, however. Its effects are very modest in increasing bone mineral density and the treatment is very expensive. Another therapy for the treatment of postmenopausal osteoporosis is the use of bis-phosphonates. These compounds were originally developed for use in Paget's disease and malignant hypercalcemia. They have been shown to inhibit bone resorption. Alendronate, one compound of this class, has been approved for the treatment of postmenopausal osteoporosis. These agents may be helpful in the treatment of osteoporosis, but these agents also have potential liabilities which include osteomalacia, extremely long half-life in bone (greater than 2 years), and possible “frozen bone syndrome,” e.g., the cessation of normal bone remodeling.\nSenile osteoporosis is similar to postmenopausal osteoporosis in that it is marked by the loss of bone mineral density and resulting increase in fracture rate, morbidity, and associated mortality. Generally, it occurs in later life, i.e., after 70 years of age. Historically, senile osteoporosis has been more common in females, but with the advent of a more elderly male population, this disease is becoming a major factor in the health of both sexes. It is not clear what, if any, role hormones such as testosterone or estrogen have in this disease, and its etiology remains obscure. Treatment of this disease has not been very satisfactory. Hormone therapy, estrogen in women and testosterone in men, has shown equivocal results; calcitonin and bis-phosphonates may be of some utility.\nThe peak mass of the skeleton at maturity is largely under genetic control. Twin studies have shown that the variance in bone mass between adult monozygotic twins is smaller than between dizygotic twins (Slemenda et al, J. Bone Miner. Res., 6:561-567 (1991); Young et al, J. Bone Miner. Res., 6:561-567 (1995); Pocock et al, J. Clin. Invest., 80:706-710 (1987); Kelly et al, J. Bone Miner. Res., 8:11-17 (1993)), and it has been estimated that up to 60% or more of the variance in skeletal mass is inherited (Krall et al, J. Bone Miner. Res., 10:S367 (1993)). Peak skeletal mass is the most powerful determinant of bone mass in elderly years (Hui et al, Ann. Int. Med., 111:355-361 (1989)), even though the rate of age-related bone loss in adult and later life is also a strong determinant (Hui et al, Osteoporosis Int., 1:30-34 (1995)). Since bone mass is the principal measurable determinant of fracture risk, the inherited peak skeletal mass achieved at maturity is an important determinant of an individual's risk of fracture later in life. Thus, study of the genetic basis of bone mass is of considerable interest in the etiology of fractures due to osteoporosis.\nRecently, a strong interest in the genetic control of peak bone mass has developed in the field of osteoporosis. The interest has focused mainly on candidate genes with suitable polymorphisms to test for association with variation in bone mass within the normal range, or has focused on examination of genes and gene loci associated with low bone mass in the range found in patients with osteoporosis. The vitamin D receptor locus (VDR) (Morrison et al, Nature, 367:284-287 (1994)), PTH gene (Howard et al, J. Clin. Endocrinol. Metab., 80:2800-2805 (1995); Johnson et al, J. Bone Miner. Res., 8:11-17 (1995); Gong et al, J. Bone Miner. Res., 10:S462 (1995)) and the estrogen receptor gene (Hosoi et al, J. Bone Miner. Res., 10:S170 (1995); Morrison et al, Nature, 367:284-287 (1994)) have figured most prominently in this work. These studies are difficult because bone mass (the phenotype) is a continuous, quantitative, polygenic trait, and is confounded by environmental factors such as nutrition, co-morbid disease, age, physical activity, and other factors. Also, this type of study design requires large numbers of subjects. In particular, the results of VDR studies to date have been confusing and contradictory (Garnero et al, J. Bone Miner. Res., 10:1283-1288 (1995); Eisman et al, J. Bone. Miner. Res., 10:1289-1293 (1995); Peacock, J. Bone Miner. Res., 10:1294-1297 (1995)). Furthermore, the work thus far has not shed much light on the mechanism(s) whereby the genetic influences might exert their effect on bone mass.\nWhile it is well known that peak bone mass is largely determined by genetic rather than environmental factors, studies to determine the gene loci (and ultimately the genes) linked to variation in bone mass are difficult and expensive. Study designs which utilize the power of linkage analysis, e.g., sib-pair or extended family, are generally more informative than simple association studies, although the latter do have value. However, genetic linkage studies involving bone mass are hampered by two major problems. The first problem is the phenotype, as discussed briefly above. Bone mass is a continuous, quantitative trait, and establishing a discrete phenotype is difficult. Each anatomical site for measurement may be influenced by several genes, many of which may be different from site to site. The second problem is the age component of the phenotype. By the time an individual can be identified as having low bone mass, there is a high probability that their parents or other members of prior generations will be deceased and therefore unavailable for study, and younger generations may not have even reached peak bone mass, making their phenotyping uncertain for genetic analysis.\nRegardless, linkage analysis can be used to find the location of a gene causing a hereditary “disorder” and does not require any knowledge of the biochemical nature of the disorder, i.e., a mutated protein that is believed to cause the disorder does not need to be known. Traditional approaches depend on assumptions concerning the disease process that might implicate a known protein as a candidate to be evaluated. The genetic localization approach using linkage analysis can be used to first find the general chromosomal region in which the defective gene is located and then to gradually reduce the size of the region in order to determine the location of the specific mutated gene as precisely as possible. After the gene itself is discovered within the candidate region, the messenger RNA and the protein are identified and, along with the DNA, are checked for mutations.\nThe genetic localization approach has practical implications since the location of the disease can be used for prenatal diagnosis even before the altered gene that causes the disease is found. Linkage analysis can enable families, even many of those that do not have a sick child, to know whether they are carriers of a disease gene and to evaluate the condition of an unborn child through molecular diagnosis. The transmission of a disease within families, then, can be used to find the defective gene. As used herein, reference to “high bone mass” (HBM) is analogous to reference to a disease state, although from a practical standpoint high bone mass can actually help a subject avoid the disease known as osteoporosis.\nLinkage analysis is possible because of the nature of inheritance of chromosomes from parents to offspring. During meiosis, the two parental homologues pair to guide their proper separation to daughter cells. While they are lined up and paired, the two homologues exchange pieces of the chromosomes, in an event called “crossing over” or “recombination.” The resulting chromosomes are chimeric, that is, they contain parts that originate from both parental homologues. The closer together two sequences are on the chromosome, the less likely that a recombination event will occur between them, and the more closely linked they are. In a linkage analysis experiment, two positions on the chromosomes are followed from one generation to the next to determine the frequency of recombination between them. In a study of an inherited disease, one of the chromosomal positions is marked by the disease gene or its normal counterpart, i.e., the inheritance of the chromosomal region can be determined by examining whether the individual displays symptoms of the disorder or not. The other position is marked by a DNA sequence that shows natural variation in the population such that the two homologues can be distinguished based on the copy of the “marker” sequence that they possess. In every family, the inheritance of the genetic marker sequence is compared to the inheritance of the disease state. If, within a family carrying an autosomal dominant disorder such as high bone mass, every affected individual carries the same form of the marker and all the unaffected individuals carry at least one different form of the marker, there is a great probability that the disease gene and the marker are located close to each other. In this way, chromosomes may be systematically checked with known markers and compared to the disease state. The data obtained from the different families is combined, and analyzed together by a computer using statistical methods. The result is information indicating the probability of linkage between the genetic marker and the disease allowing different distances between them. A positive result can mean that the disease is very close to the marker, while a negative result indicates that it is far away on that chromosome, or on an entirely different chromosome.\nLinkage analysis is performed by typing all members of the affected family at a given marker locus and evaluating the co-inheritance of a particular disease state with the marker probe, thereby determining how often the two of them are co-inherited. The recombination frequency can be used as a measure of the genetic distance between two gene loci. A recombination frequency of 1% is equivalent to 1 map unit, or 1 centiMorgan (cM), which is roughly equivalent to 1,000 kb of DNA. This relationship holds up to frequencies of about 20% or 20 cM.\nThe entire human genome is 3,300 cM long. In order to find an unknown disease gene within 5-10 cM of a marker locus, the whole human genome can be searched with roughly 330 informative marker loci spaced at approximately 10 cM intervals (Botstein et al, Am. J. Hum. Genet., 32:314-331 (1980)). The reliability of linkage results is established by using a number of statistical methods. The method most commonly used for the analysis of linkage in humans is the LOD score method (Morton, Prog. Clin. Biol. Res., 147:245-265 (1984), Morton et al, Am. J. Hum. Genet., 38:868-883 (1986)) which was incorporated into the computer program LIPED by Ott, Am. J. Hum. Genet., 28:528-529 (1976). LOD scores are the logarithm of the ratio of the likelihood that two loci are linked at a given distance to that they are not linked (>50 cM apart). The advantage of using logarithmic values is that they can be summed among families with the same disease. This becomes necessary given the relatively small size of human families.\nBy convention, a total LOD score greater than +3.0 (that is, odds of linkage at the specified recombination frequency being 1000 times greater than odds of no linkage) is considered to be significant evidence for linkage at that particular recombination frequency. A total LOD score of less than −2.0 (that is, odds of no linkage being 100 times greater than odds of linkage at the specified frequency) is considered to be strong evidence that the two loci under consideration are not linked at that particular recombination frequency. Until recently, most linkage analyses have been performed on the basis of two-point data, which is the relationship between the disorder under consideration and a particular genetic marker. However, as a result of the rapid advances in mapping the human genome over the last few years, and concomitant improvements in computer methodology, it has become feasible to carry out linkage analyses using multi-point data. Multi-point analysis provide a simultaneous analysis of linkage between the disease and several linked genetic markers, when the recombination distance among the markers is known.\nMulti-point analysis is advantageous for two reasons. First, the informativeness of the pedigree is usually increased. Each pedigree has a certain amount of potential information, dependent on the number of parents heterozygous for the marker loci and the number of affected individuals in the family. However, few markers are sufficiently polymorphic as to be informative in all those individuals. If multiple markers are considered simultaneously, then the probability of an individual being heterozygous for at least one of the markers is greatly increased. Second, an indication of the position of the disease gene among the markers may be determined. This allows identification of flanking markers, and thus eventually allows isolation of a small region in which the disease gene resides. Lathrop et al, Proc. Natl. Acad. Sci. USA, 81:3443-3446 (1984) have written the most widely used computer package, LINKAGE, for multi-point analysis.\nThere is a need in the art for identifying the gene associated with a high bone mass phenotype. The present invention is directed to this, as well as other, important ends."}
{"text": "Vehicular interior lighting, e.g., dome lights, courtesy lights, glove box lights, foot-well lights, and other roof and side panel interior lighting, has generally been realized through the application of conventional, incandescent lighting sources for the past few decades. Although the light produced from incandescent sources is fairly uniform, it is generally accompanied by relatively high heat levels and low intensity compared to more modern light sources (e.g., light-emitting diode, fluorescent, etc.). Further, incandescent light sources have relatively high energy usage compared to more modern light sources.\nIn the past decade, light emitting diode (LED) sources have been employed in various lighting applications, including vehicular applications. LED sources possess many advantages over conventional lighting sources including the ability to transmit high intensity light patterns and generate low heat upon light transmission with very low energy usage. One drawback of LED sources is that the light produced from them is directional and can result in high-contrast illumination, shadowing and other undesirable effects associated with highly directional light. Consequently, in the context of vehicular interior lighting, LED sources have been difficult to use given that the directional nature of the light output from these sources has led to high contrast within the close quarters of the vehicle interior.\nSome success has been realized in overcoming the directional nature of the light output from LED sources by coupling certain optics to them. In particular, near-field lens elements, collimators, light-diffusers and reflectors have been coupled to LED sources to shape their light output into uniform patterns approximating those produced by incandescent sources. Unfortunately, many of these solutions require additional packaging (e.g., larger overall lighting source dimensions) that can significantly add cost; consequently, LED sources are not prevalent in vehicular interiors, particularly as replacements to low-cost courtesy and dome lights\nConsumers are also increasingly demanding access to decorative lighting effects within vehicles, residences, offices and other interiors. In vehicles, decorative lighting patterns that supplement the natural light transmitted through side windows, windshields, and rear windows is often provided through lighting assemblies (e.g., fiber optics) that are independent of sun roofs, moon roofs and other non-traditional windows. Many consumers, however, find significant drawbacks associated with these separate lighting assemblies, including the loss of interior space associated with the elements needed for these decorative lighting access points, maintenance concerns, and added cost.\nAccordingly, there is a need for vehicular interior lighting solutions that can employ LED sources in a low-cost fashion, with minimal packaging constraints and at low manufacturing costs. Further, as LED sources continue to be integrated within the vehicle industry, there is a desire by many consumers for unique lighting aesthetics that cannot be achieved through conventional sources."}
{"text": "The disclosure relates to a hydraulic axial piston machine which has a housing, a drive shaft which is rotatably mounted in the housing, a swashplate, which for the purpose of varying its inclination relative to the axis of the drive shaft can be pivoted about a pivot axis, and an actuating piston, which extends with its longitudinal axis spaced apart from a central plane, which is perpendicular to the pivot axis and extends through the axis of the drive shaft, of the swashplate and at least approximately parallel to the axis of the drive shaft. The actuating piston is situated, so to speak, very far outward in a corner of the housing. Said actuating piston, at a first end, engages on the swashplate for the purpose of adjusting the latter and, at a second end, delimits an actuating chamber into which control fluid flows for the purpose of pivoting the swashplate in one direction and out of which control fluid can be displaced in the event of a pivoting movement of the swashplate in the other direction. On the actuating piston there is arranged an elongate feedback element by means of which the position of the actuating piston and thus the oblique position of the swashplate are input to control a regulating valve.\nIn a hydraulic axial piston machine of said type known from DE 10 2007 022 569 A1, the longitudinal axis of the actuating piston and the longitudinal axis of the feedback element which projects perpendicularly from the actuating piston span a plane which is perpendicular to the pivot axis of the swashplate.\nA feedback element of said type is provided in particular if the axial piston machine is to be adjusted in a torque-regulated manner, or proportionally to an input signal. In the case of torque regulation, the feedback element is also provided with a small piston which is subjected to the working pressure and which, depending on the position of the actuating piston and thus of the swashplate, engages on a lever at a different distance from an axis of rotation and exerts a torque on said lever. The valve piston of a regulating valve is supported counter to said torque on the same arm, or on a second arm, of the lever at a fixed distance from the axis of rotation, said valve piston being subjected to a constant or remote-controlled variable force which seeks to increase the swept volume. The swept volume of the axial piston machine is then set in each case such that torque equilibrium prevails at the lever.\nIn the case of a proportional adjustment of the swept volume, the feedback element varies the preload of a spring which exerts load on a valve piston of the regulating valve, said valve piston being acted on counter to the spring by an input force generated predominantly by an electromagnet or a hydraulic pressure. Depending on the magnitude of the input force, the spring force and thus the position of the actuating piston and thus of the swashplate must vary such that, when the valve piston is in the zero position, the spring force and input force maintain the equilibrium.\nA high degree of flexibility with regard to the arrangement of the regulating valve should be possible.\nThis is achieved in that the feedback element is positioned within its range of movement in each case such that its longitudinal axis and the longitudinal axis of the actuating piston span a plane which differs from a plane perpendicular to the pivot axis of the swashplate. The positioning of the feedback element is determined for example by a guide in the housing or on the regulating valve or by a particular arrangement on the actuating piston if the latter is not rotatable about its longitudinal axis.\nAccording to the disclosure, it is possible for a regulating valve to be arranged at different locations on the outside of the housing. It is necessary in each case merely to provide a corresponding opening on the housing. The actuating piston and feedback element and also the regulating valve can always be identical.\nAdvantageous embodiments of a hydraulic axial piston machine according to the disclosure emerge from the description below."}
{"text": "1. Technical Field\nThe invention relates to a retaining wall, in particular a retaining wall inclined against a soil backfill, composed of individual shaped blocks, the shaped blocks being arranged one above the other in layers and the shaped blocks of adjacent layers each engaging in a form-fitting manner with one another in the region of upper sides and lower sides of the shaped blocks.\n2. Prior Art\nSuch retaining walls are known from practice in numerous variants. It is known here both to bond together the layers of shaped blocks with mortar or to construct them as a so-called gravity retaining wall without mortar. It is also known to couple the individual shaped blocks of the retaining wall by means of connecting pieces.\nA retaining wall of the type mentioned initially is known for example from German utility model 90 15 718, originating from the applicant. The form-fitting connection of the shaped blocks in the region of the upper sides and lower sides produces a bond within the layers. This results in a wall element which, due to the bond of the shaped blocks, can be viewed statically as a plate element. The same applies when the shaped blocks are mortared.\nA problem with the above-described retaining wall is that the plate element can fail as a whole under varying loading, with the result that the retaining wall is damaged or even destroyed. Such loading variations can be produced, for example, by alternating loading by the soil backfill, in particular by alternating applied loads or alternating moisture states."}
{"text": "1. Field of the Invention\nThe present invention relates to a damper device located on a liquid supply path connecting a liquid supply that supplies a liquid and a liquid injector that injects the liquid, and relates to a liquid supply system including such a damper device, and an inkjet recording device.\n2. Description of the Related Art\nAn inkjet recording device for industrial use or the like adopts a structure in which a large capacity ink cartridge is located away from a carriage having an ink injection head mounted thereon, namely, an off-carriage system. In a recording device of the off-carriage system, a large dynamic pressure fluctuation is generated by a movement of a carriage. If the large dynamic pressure fluctuation is applied to the ink injection head, ink injection may be destabilized or ink may be unexpectedly dropped from the injection head to stain the printed item.\nIn such a situation, it has been conventionally studied to suppress the ink pressure fluctuation in the ink injection head. For example, Japanese Laid-Open Patent Publication No. 2005-059274 discloses a structure in which an ink injection head and a damper device are mounted on a carriage.\nThe damper device described in Japanese Laid-Open Patent Publication No. 2005-059274 includes a case main body forming an ink storage chamber, a flexible pressure sensitive film acting as one surface of the ink storage chamber, and an elastic member urging the pressure sensitive film externally away from the ink storage chamber. At substantially the center of the pressure sensitive film, a reinforcing member is bonded. In this damper device, ink is supplied at a predetermined pressure to the ink injection head based on a flexible deformation of the pressure sensitive film.\nHowever, in the damper device described in Japanese Laid-Open Patent Publication No. 2005-059274, in the case where, for example, the reinforcing member is bonded to the pressure sensitive film in a certain manner, the pressure sensitive film may not be flexibly deformed in a proper manner and thus the ink may be supplied unstably. In the case where, for example, substantially the entire surface of the reinforcing member is bonded to the pressure sensitive film, the movable area of the pressure sensitive film may be narrowed. In the case where air is confined at a bonding surface when the reinforcing member is bonded to the pressure sensitive film, the air may act as a resistance and as a result, the pressure sensitive film may be flexibly deformed irregularly."}
{"text": "Absorbent articles such as diapers, training pants, incontinence briefs, and the like are well known for their convenience in absorbing and containing body exudates. Typical absorbent article construction involves various structures having defined properties aiding the fit or functionality of the article. However, it is advantageous to reduce the number of structures and amount of material present in absorbent articles. For disposable absorbent articles, fewer structures and less material results in reduced disposable mass. Fewer structures may also simplify the processing of the absorbent article. Fewer structures may also reduce the total cost of the absorbent article. One way to reduce the number of structures and amount of material present in an absorbent article is by combining the function of multiple structures into a single structure. The core support and the barrier leg cuff are two such structures that may be combined.\nThe core support is a substrate upon which the absorbent core is disposed. The core support typically comprises an air permeable material such as a nonwoven web, a cellulosic tissue, an apertured film, or other like materials. The core support may assist in maintaining the integrity of the absorbent core. In particular, the cores of modern absorbent articles comprise relatively high percentages of superabsorbent polymer (SAP) particles. Prior to liquid insult, the SAP particles often resemble grains of sand. The core support may be used to confine the SAP particles within the core and to prevent the SAP particles from becoming distributed throughout the absorbent article. Additionally, the core support may be a necessary component of core formation. The core support may serve as the formation layer upon which the components of the core are placed during the process of forming the core. Typical core components include SAP particles and cellulosic fluff. During the formation of the core, a vacuum is often drawn through the core support so that core components disposed on the core support remain immobilized. This type of vacuum formation necessitates a core support that has a relatively high degree of air permeability, an air flow measured in, for example, m3/m2/min, such that a vacuum may be drawn through the core support.\nBarrier leg cuffs (e.g., also referred to as inner cuffs, inner leg cuffs, leg gussets, standing leg cuffs, barrier cuffs) are physical barriers which inhibit loose fecal material, urine, or liquids from escaping the article. The barrier cuffs restrain the free movement of exudates and provide a structure to contain the exudates within the diaper. Typical barrier cuffs include a pair of flaps disposed longitudinally on the article and running at least through the crotch region of the article. The barrier cuffs are laterally spaced so as to allow for the receipt of body exudates into the absorbent article. Barrier cuffs typically comprise an elastic member associated with the barrier cuff that allows the barrier cuff to stand up and serve as a physical barrier to exudate leakage or run-off from the body-facing surface of the article. The flaps may comprise a liquid impermeable material. Suitable liquid impermeable materials include materials that are substantially or fully liquid impermeable or may be treated to become more liquid impermeable such as woven webs, nonwoven webs, films, and other like materials. Another consideration for flap construction is breathability. During wear, a portion of the barrier leg cuff is in contact with the wearer. A barrier leg cuff having a breathable flap is desirable to prevent occlusion and over hydration of the skin. The flap may be constructed from a material that has some degree of air and vapor permeability.\nCombining the core support and the barrier leg cuff into a single structure requires that the structure perform different functions at different locations in the diaper. For example, the single structure may require, in some regions, the air permeability and SAP particle retention properties of the core support and may require, in other locations, the liquid impermeability and air/vapor permeability of the barrier leg cuff. While integrating barrier leg cuffs with a substrate underlying the core has been tried in absorbent articles (e.g., U.S. Pat. No. 5,643,239), previous executions have failed to appreciate the multiple functions that the integrated structure must exhibit. Particularly, previous applications have failed to appreciate the dynamic between high air permeability in certain zones (e.g., for core formation) and liquid impermeability in other zones (e.g., for barrier protection). As a result, previous executions disclose an integrated structure that may compromise one of the functions that the structure was intended to perform (e.g., requisite air permeability versus liquid impermeability)."}
{"text": "The replacement of destroyed or damaged human joints is one of the great achievements of twentieth century orthopaedic surgery. However, total joint prostheses, composed of various combinations of metal, ceramic, and polymeric components, continue to suffer from distressingly limited service lives. For example, the current generation of high-load bearing prostheses for the hip and knee have a typical lifetime on the order of 6-12 years. Generally, failed implants can be replaced once or twice, which means that current technology provides a solution for--at most--about 25 years.\nWith human life expectancies steadily increasing, there is a driving need to increase significantly the effective lifetime of a single implant. One of the problems encountered in designing such prostheses is the difficulty of finding materials which are both biocompatible and also durable enough to replace a human joint. In use, a human joint is exposed to substantial, repetitive loads and frictional stresses.\nAlthough the geometric details may vary, a natural human hip, knee, or shoulder joint generally includes: (a) a more-or-less spherical ball; (b) an attachment to a long bone; and, (c) a hemispherical socket (the \"acetabular cup\") in a contiguous bony structure which retains the spherical ball so that the long bone may pivot and articulate. In a healthy joint, nature minimizes the friction between the joint components and prevents bone-on-bone wear and destruction by using mating porous cartilaginous layers that provide \"squeeze-film\" synovial fluid lubrication. This lubrication results in a low coefficient of friction on the order of 0.02.\nWhen a human joint has been destroyed or damaged by disease or injury, surgical replacement (arthroplasty) normally is required. A total joint replacement includes components that simulate a natural human joint, typically: (a) a more-or-less spherical ceramic or metal ball, often made of cobalt-chromium alloy; (b) attached to a \"stem,\" which generally is implanted into the core of the adjacent long bone; and (c) a hemispherical socket which takes the place of the acetabular cup and retains the spherical ball. This hemispherical socket typically is a metal cup affixed into the joint socket by mechanical attachments and \"lined\" with UHMWPE so that the ball can rotate within the socket, and so that the stem, via the ball, can pivot and articulate.\nOne of the difficulties in constructing any device for implantation into the human body is the need to avoid an adverse immune response. The possibility of an adverse immune response is reduced when certain synthetic materials are used. Cobalt-chromium alloy, titanium, and UHMWPE are examples of such synthetic materials. Unfortunately, the use of UHMWPE in the bearing of a total joint replacement may be the cause of at least one type of failure of such devices.\nThree basic problems may cause a total joint replacement to fail or to have a limited service life. The first problem, which manifests itself at the bone-stem interface, is not the focus of the present application. Because the elastic modulus of the stem greatly exceeds that of the bone, flexural loading caused by walking creates local cyclic stress concentrations due to the non-compliance of the stem. These stresses can be intense and even severe enough to cause death of local bone cells. If this occurs, pockets of non-support are created, and the stem may loosen or fail.\nThe two other basic problems, which are coupled, are the subject of the present application. One of these problems, known as ball-cup friction and wear, results from frictional wear between the hemispherical bearing (which is \"lined\" with UHMWPE) and the polished spherical ceramic or metal ball attached to the stem. The other problem, known as sub-surface fatigue, results from brittleness of the UHMWPE bearing and the resulting tendency of the UHMWPE bearing to fail under reciprocating applied loads.\nFor many years, the acetabular cup in joint implants has been \"lined\" with UHMWPE, or other materials, in order to decrease the coefficient of friction of the socket or bearing. Unfortunately, clinical experience has shown that, at least when UHMWPE is used to line the bearing, either the surface of the UHMWPE \"bearing\" and/or the surface of the metal/ceramic ball ultimately is destroyed by friction-induced wear. Alternately, the acetabular cup loosens after a period of use, greatly increasing ball-cup friction and wear.\nSome insight into the cause of failure due to ball-cup friction and wear has been gleaned from histological studies of the surrounding tissue. These histological studies show that the surrounding distressed tissue typically contains extremely small particles of UHMWPE which range from sub-micrometers to a few micrometers in size. Larger particles of UHMWPE appear to be tolerated by the body, as is the solid bulk of the UHMWPE bearing. However, the body apparently does not tolerate smaller particles of UHMWPE. In fact, these small particles of UHMWPE cause powerful histiocytic reactions by which the body unsuccessfully attempts to eliminate the foreign material. Agents released in this process attack the neighboring bone to cause \"wear debris-induced osteolysis\" which, in turn, leads to a loss of fixation and loosening of the prosthesis due to \"remodeling\" of the bone.\nThe first step in the generation of the small particles of UHMWPE appears to be the formation of a very thin layer of polyethylene between the spherical ball and the UHMWPE lining of the bearing. This thin film of polyethylene adheres to the \"ball\" and serves as a soft, shearable, solid lubricant composed of millions of submicrometer particles. Adhesive wear between the ball and the bearing produces strong, adhesive junctions on the ball. When exposed to further friction, fibrils of the polymer shear off of these adhesive junctions and are drawn into slender connecting ligaments, eventually producing ligament rupture.\nThis ligament rupture apparently produces the lubricous, extremely small particles of UHMWPE which eventually migrate to the bone-acetabular cup bond line. The reason for migration of these particles into the \"crevice\" between the ball and the cup are the microcurrents that are generated in the synovial fluid by joint motion. Once a sufficient number of small particles enter the bone-cup crevice, the bone tissue begins to degrade and the joint replacement eventually loosens and fails.\nOne way to reduce friction between the metal and UHMWPE components would be to coat one or both of the components with diamond-like carbon (DLC), which is chemically inert, biocompatible, and is known to have a low coefficient of friction. Unfortunately, the very properties of DLC that make it a desirable coating for parts that will be frictionally engaged make it difficult to achieve strong adhesion of the DLC coating to the substrate, particularly where deposition temperatures must be low. This limited adhesion problem can be exacerbated by very high compressive stress, such as that found in a plasma-deposited DLC (up to 8 GPa). Therefore, some have concluded that DLC--or at least plasma-deposited DLC--cannot be used in orthopaedic applications.\nEnergetic ion beam-associated DLC has a far lower residual stress than plasma-deposited DLC, and is a better candidate for a high integrity DLC. The substrate material to which all forms of carbon adhere most successfully is silicon. This is because strong covalent Si--C bonds are easily formed between the coating and the silicon substrate. Some have attempted to improve the adhesion of DLC to other materials, such as metal alloys, by forming an interposed silicon bond-coat to which the DLC will adhere more strongly.\nUnfortunately, this simple approach does not result in adhesion that survives in applications, such as orthopaedic applications, where the DLC coating is subjected to substantial friction and stress. The simple formation of a silicon bond-coat on a metal alloy appears to create another relatively weak interface between the silicon and the metal or alloy.\nTherefore a method is needed by which a DLC coating can be strongly adhered to a metal surface, and by which the shearing of polymer fibrils from an UHMWPE component can be prevented. The method would be most efficient if it rendered the UHMWPE compound less brittle so that sub-surface fatigue failure was reduced."}
{"text": "This invention relates to refuse collection vehicles.\nMore particularly, the present invention relates to the fabrication of refuse collection semi-trailers.\nIn a further and more specific aspect, the present invention concerns a method of converting a refuse collection vehicle into a semi-trailer.\nThe collection and removal of refuse, the solid waste of a community, is a major municipal problem. Traditionally, residential refuse, including garbage, trash, and other waste materials, is collected and stored in containers having a capacity of 10 to 30 gallons, although these figures can vary depending on the community. Commercial or industrial refuse generators accumulate refuse in larger heavier containers.\nConventionally, refuse containers are emptied into a refuse collection vehicle which transports the collected refuse to a disposal site such as a land fill, dump, incinerator, etc. The refuse collection vehicle includes a vehicle chassis with a cab mounted at one end and a refuse collection body mounted on the chassis behind the cab. A hopper is typically included, into which refuse is deposited for transfer into the refuse collection body.\nHoppers can be located at the front of the body or the rear. Refuse container handling devices are often used to deposit refuse from the containers into the hoppers. Front loaders and side loaders typically deposit refuse into a hopper located at the front of the vehicle. Rear loading vehicle are often manually loaded but may also be loaded by various mechanical devices.\nRefuse collection vehicles typically must travel quite a distance to reach a disposal site. As sites closer to city centers become full or operate at capacity, new sites are opened further away from the areas of refuse collection. Because landfill regulations increase the cost of preparation and closing, the trend is for fewer disposal sites at greater distances apart. This increases the travel time needed to dispose of refuse. Conventional refuse collection vehicles are limited in the amount of refuse they can carry by the weight supported by the rear axles. This is particularly troublesome in rear loading vehicles in which the heavy rear loading hopper mechanism is positioned behind the rear axles. Thus a larger proportion of the weight is supported by the rear axle reducing load capacities. A longer vehicle will carry a larger load, but once the weight limit on the rear axle is reached, more space becomes irrelevant. Additionally, a longer vehicle will become unwieldy and difficult to maneuver.\nThese problems have been solved by the use of semi-trailer refuse collection vehicles as disclosed in U.S. Pat. No. 5,551,824, entitled xe2x80x9cArticulated Refuse Collection Apparatusxe2x80x9d filed in July 1994. This reference teaches the use of semi-trailers to provide larger carrying capacity as well as greater maneuverability. However, manufacturing of these semi-trailers can be very expensive. Manufacture of semi-trailers require retooling of equipment and establishment of new production lines. Many municipalities and companies are not interested in renovating their fleets of refuse collection vehicles due to these costs involved.\nIt would be highly advantageous, therefore, to remedy the foregoing and other deficiencies inherent in the prior art.\nAccordingly, it is an object the present invention to provide a method of converting a refuse collection body into a semi-trailer.\nBriefly, to achieve the desired objects of the present invention in accordance with a preferred embodiment thereof, provided is a method of converting a refuse collection body into a semi-trailer. A refuse collection body having a front end and a rear end, and a frame extending along a bottom of the body from adjacent the front end to adjacent the rear end is provided. A king pin is mounted to the frame proximate the front end for attachment to a fifth wheel of a tractor and an axle assembly is mounted to the frame proximate the rear end.\nIn a preferred embodiment, the frame is shaped proximate the front end to form a contour of the frame substantially matching a surface of the fifth wheel of a tractor. Shaping the frame includes removing portions thereof and affixing cover portions. The frame includes spaced apart, elongated hollow members each having a bottom wall and sidewalls, and the step of removing portions of the frame includes removing a portion of the bottom wall and portions of the sidewalls. Cover portions are then affixed to the remaining portions of the sidewalls forming another bottom wall."}
{"text": "Up to 170 million people worldwide are chronically infected with HCV, with many at significant risk for cirrhosis, liver failure and hepatocellular carcinoma. The World Health Organization estimates an annual increase in the global burden by 2 million new infections (Shepard et al. (2005) Lancet Infect Dis 5:558-567). Efforts at vaccine development have yet to succeed despite two decades of work. A key step towards reaching this goal is to identify and exploit relevant mechanisms of immune protection. Although there are no clearly established in vitro correlates of protective immunity, multiple lines of evidence suggest that CD4+ and CD8+ T cell responses, while critical for controlling acute HCV infection, are inadequate for prevention of long-term persistence in most infected individuals (Bowen et al. (2005) Nature 436:946-952). Emerging evidence supports the importance of virus-neutralizing (Vn) antibodies, and the ability of B cell responses to modify the course of infection.\nThe development of in vitro cell culture models, based on HCV retroviral pseudotype particles expressing E1E2 (HCVpp) and infectious cell culture-derived HCV virions (HCVcc), has facilitated the measurement of antibody-mediated virus neutralization, and thus to evaluate the impact of antibody in the control of infection (Bartosch et al. (2003) J Exp Med 197:633-642, Cai et al. (2005) J Virol 79:13963-13973, Hsu et al. (2003) Proc Natl Acad Sci USA 100:7271-7276, Lindenbach et al. (2005) Science 309:623-626, Wakita et al. (2005) Nat Med 11:791-796, Zhong et al. (2005) Proc Natl Acad Sci USA 102:9294-9299). Chimpanzee studies have shown that protection from an infectious HCV inoculum with HCV-specific IgG is correlated with antibody titers blocking infection of target cells with HCVpp (Bartosch et al. (2003) Proc Natl Acad Sci USA 100:14199-14204). Control of virus infection and the Vn antibody response measured via HCVpp have been associated in single source outbreaks of acute HCV infection (Lavillette et al. (2005) J Virol 79:6023-6034, Pestka et al. (2007) C. Proc Natl Acad Sci USA 104:6025-6030), and confirmed in a study of active injection drug users (Osburn et al. (2010) Gastroenterology 138:315-324). While only 25% of subjects in this study cleared primary HCV infection, 83% cleared re-infection, and clearance in some subjects was associated with cross-reactive Vn antibodies.\nIn addition, antibodies to HCV E2 prevent infection in a human liver-mouse chimeric model (Law et al. (2008) Nature Medicine 14:25-27, Meuleman et al. (2008) Hepatology 48:1761-1768). Finally, an immunocompetent humanized mouse model for HCV exhibited a robust antibody response to a recombinant vaccinia virus expressing HCV proteins that protected against a heterologous infectious HCV challenge in some of the animals, which was correlated with the serum level of antibodies to E2 (Dorner et al. (2011) Nature 474:208-211).\nA significant challenge for vaccine development is defining conserved epitopes that i) are capable of eliciting protective antibodies in this highly diverse virus, and ii) are resistant to development of escape mutants. Treatment of HCV and the development of vaccines that broadly protect against highly diverse HCV genotypes and subtypes are of interest in the field. The present invention addresses this issue."}
{"text": "Arrays of hard disk drives (HDDs)connected to host computer systems are commonly used for computer data storage. Disk drive arrays provide large storage capacities and high reliability at a low cost.\nU.S. Pat. No. 4,870,643 teaches an array of disk drives where a set or stripe of data blocks and a parity block for the set are stored on separate disk drives, with no one disk drive having all the parity for all of the sets. In the event a disk drive fails, each block of data on the unavailable disks can be reconstructed using the remaining blocks in the set. Whenever a data block is updated with new data the associated parity is also updated. This is often referred to as a RAID level 5 system.\nIn Patterson et al., \"A case for Redundant Arrays of Inexpensive Disks (RAID)\", A. C. Sigmod Conference, Chicago, Ill., Jun. 1-3, 1988, pp. 109-116, five levels of RAID were defined. In each of the RAID levels, an array of disk drives includes redundant data which can be used to reconstruct data in the event one of the drives fails. RAID level 1 refers to data mirroring where a duplicate copy of the data on one disk drive is stored on a separate disk drive. RAID level 2 uses Hamming codes to provide error detection and correction for a data set. In a RAID level 3 system, a block of data is divided into N portions for storage on N disk drives. The portions of the data block are exclusive ORed (XOR) to produce parity which is written to a separate N+1 disk drive. In a RAID level 4 system, blocks of data are stored on separate disk drives with the parity (XOR) information for a set of blocks of data being stored on a separate dedicated disk drive. The set of blocks are referred to as a stripe of data blocks. Typically, the data blocks and the parity block of a set or stripe are written to the same logical block address on respective disk drives. A RAID level 5 system is similar to RAID level 4 system except that the parity is distributed among all of the disk drives. In RAID level 5 arrays, disk drives operate independently of each other, so that multiple read and write operations can access separate blocks of data at the same time.\nThe RAID advisory board in its RAID book, \"A Source Book for Disk Array Technology\", 5th Edition, recognizes these five levels of RAID and defines a sixth level of RAID. In RAID-6, a second independent parity block is provided for each set, so that there are N+2 member disks. This allows data on two failed disk drives to be reconstructed.\nIn all RAID systems, when data is written (updated), the corresponding redundant data needs to be updated as well. In RAID levels 4 through 6 systems, where the data blocks are independently accessed on separate disk drives, during a write operation, parity can be updated by XORing the old data, the new data and the old parity. Alternatively, a set of data blocks can be written together with its parity to the set of disk drives.\nIn most prior art RAID storage systems, a RAID controller manages the storage and retrieval of the data on the array of disk drives and the parity generation. Array management software running on one host system or in a storage sub-system manages the storage and retrieval of data to and from the storage devices. Application programs running on the host system provide a logical disk address for writing and retrieving data. The array management software translates the address and reads the contents of the requested block into the host memory. When modifying data, the array management software on the host or a storage controller reads the contents of the block to be modified and the corresponding parity block, and then calculates an exclusive OR (XOR) on the old data, old parity, and the new data. The array management software then writes the new parity back to its parity location and writes the new data to its prior location. A RAID controller can be implemented purely in software or a combination of microcode and hardware. The controller and parity generator can reside in the host computer system or in a separate storage subsystem.\nMore recently, array management systems have been designed where array management functions and parity generation are performed within the storage devices rather than using a separate controller. More specifically, disk drives have been designed with the XOR engine incorporated into the drive. The disk drives rely on peer-to-peer communication over an interface, such as the Small Computer Standard Interface (SCSI), to manage the implementation of the parity updates and reconstruction.\nPerforming XOR operations in the disk drive can result in reduced data transfers across the interconnections between the disk drives and the host system.\nIn a host based or sub-system based array management system, when a RAID level 5 write operation is performed, there are four data transfers that are involved and the array controller executes two XOR operations. These four transfers are: (1) transferring the old data to the controller (2) transferring the old parity to the controller (3) transferring the new data from the controller, and (4) transferring the new parity from the controller.\nIn a RAID system where the drives perform the XOR operations, data can pass directly from drive to drive for the XOR operations. This greatly reduces the amount of work that a separate host computer or controller has to perform and reduces the amount of data transfers over the interfaces between the drives and the controller. For a write operation (which includes updating parity) the number of data transfers is reduced from four to two. The two data transfers are (1) transfering the new data to the disk; and (2) transfering the XOR difference to the disk where the corresponding (old) parity is stored. In such an array, a new data block is transferred from the host to the disk drive where the old data block is stored. The disk drive accepts the new data, reads the corresponding old data from its disk, and performs an XOR operation to determine the XOR difference between the old data and the new data. The disk drive then acts as an initiator to transfer the old data/new data XOR difference to the disk drive that has the corresponding parity for the data block being updated. The disk drive that has the parity, accepts the XOR difference and performs an XOR operation between the XOR difference and the old parity (read from its disk) to produce the new parity. The new parity is then written back to its disk. By performing the XOR operation in the disk drives, there is also no longer a need for the array controller to perform XOR operations.\nThe ANSI Standards Committee established commands for implementing RAID functions on a set of disk drives having XOR function capabilities. See: \"XOR commands on SCSI Disk Drives\" X3T/1/96-IIIR2. Such a system is referred to as an XOR-on-the drive system. The current proposed ANSI standard for \"XOR commands\" on SCSI disk drives includes read (READ), data update (XDWRITE), parity update (XPWRITE), and reconstruction (REGENERATE and REBUILD) commands. The XDWRITE, REGENERATE and REBUILD commands are executed by drives acting as temporary initiators using peer-to-peer communication. There have also been enhancements to the current proposed SCSI standard. Commonly owned patent application, Hodges, \"A System and Method for Distributing Parity in an Array Operation\" Ser. No. 08/396,046, teaches a system where the drives store information on the RAID configuration which otherwise would be provided by the host. These enhancements use commands similar to the proposed standard.\nThe READ command is the most common proposed ANSI command. A host requesting a specific data block issues a READ command to the appropriate disk drive to read that block of data.\nUpdate operations involve writing data on one drive of the array and updating corresponding parity information on a second drive. There are two commands in the proposed ANSI standard that are used to accomplish an update, XDWRITE and XPWRITE. The host issues an XDWRITE command and sends new data to the disk drive where the old data is stored. This disk drive then acts as an initiator. The initiator calculates a parity difference between the old and new data by XORing the new data with the old data read from its disk. The initiator issues an XPWRITE command sending the calculated parity difference between the new and old data to the drive containing the corresponding parity for the data. The parity drive XOR's the parity difference with the old parity, read from its disk, in order to produce the new parity, which is then written back to its disk.\nReconstruction operations involve reading data from multiple drives of the redundant array and performing an exclusive OR operation to recover data. There are two commands in the proposed ANSI standard relating to the reconstruction of data, REGENERATE and REBUILD.\nThe REGENERATE command is used in place of a READ command when a data drive in the array has malfunctioned. A host computer sends a known good drive the REGENERATE command with the addresses of the source drives (i.e. the drives storing the data and parity blocks of the unavailable block's stripes), the corresponding logical block addresses of the data and parity blocks, and the number of blocks to be reconstructed. The known good drive takes on the role of the initiator and sends READ commands to all of the other drives which have the corresponding data blocks and parity blocks in the same set or stripe as the data block to be regenerated, and also reads the data from its own disk. The blocks from all the drives are then sent back to the initiator drive where the blocks are exclusive ORed, and the result is then sent to the host.\nThe host sends a REBUILD command to a replacement drive to be rebuilt, which acts as an initiator. The drive has the addresses for the other source drives as well as the logical block addresses and the number of blocks to be rebuilt. The initiator issues READ commands to all of the source drives having data blocks or parity blocks in the parity stripe of the block to reconstructed. When the drive receives all of these blocks, it exclusive OR's the blocks and writes the result of the exclusive OR operation to its disk.\nReconstruction requires that data from all devices be read at the same logical point in time so that the data blocks of a stripe and the corresponding parity block are consistent. If reconstruction and update operations are performed on the same data blocks without specific coordination there is a possibility of incorrect reconstruction. To ensure consistency no writes may be permitted to data blocks of the stripe on any drive until all of the stripe data and parity blocks have been read from all devices for the data to be reconstructed. This is readily accomplished by a single RAID controller. A single controller can serialize accesses to its attached DASD as needed. It is more difficult in a distributed RAID controller environment such as presented by XOR-in-drive configurations. The present ANSI definitions leave the serialization and coordination of the operations to the host system. The problem of incorrect reconstruction arises when the reconstructed and updated operations are performed by two independent initiators, each initiator being unaware of what the other initiator is doing.\nFIG. 1 shows an example where a disk drive controlled array tries to implement conflicting command executions for blocks in the same parity stripe. An application program running on a host 10 update writes a data block B3, stored on drive D4, while the same or a different host reads a data block B4, from the same parity stripe as block B3, that is stored on drive D5. Since the drive D5 has failed, the data block B3 will need to be reconstructed.\nIn order to reconstruct data block B4 stored on the failed drive, the host issues the command REGENERATE to drive D1. Drive D1 acts as an initiator for the regeneration function using the corresponding data and parity blocks from B4's parity stripe stored on drives D1, D2, D3 and D4.\nShortly after drive D1 receives the REGENERATE command, drive D4 receives an XDWRITE command to update block B3. The update write of block B3 requires the updating of the corresponding parity block for the stripe which is stored on drive D2.\nAs shown in FIG. 1, without the proper coordination of the execution of the commands, the regeneration operation incorrectly uses the old parity block with the new data block for a parity stripe.\nDrive D1 issues commands to drives D2, D3 and D4 to read blocks P, B2, and B3 12 and reads the corresponding data block B1 from its own disk 14. When all of the data and parity blocks are received, drive D1 XORs the data and parity blocks 16.\nWhen drive D3 receives the READ command issued from drive D1 (as part of the regeneration), there are no other commands waiting to execute, so drive D3 executes the READ command from D1 and sends the read data to drive D1 18.\nMeanwhile, drive D4 executes the XDWRITE command for block B3 and acting as a second separate initiator issues an XPWRITE command to drive D2 (where the corresponding parity P is stored) 20. Drive D4 enqueues the READ command issued from drive D1 (for the regeneration) 22. After drive D4 completes the execution of the XDWRITE command by reading the old data, sending the parity difference to drive D2, and writing the updated data block to its disk 24, drive D4 then executes the READ command P from its queue (issued from drive D1) 25. At that point, drive D4 is reading the updated data block.\nWhen drive D2 receives the READ command from drive Dl (as part of the REGENERATE function) 26 and executes it, drive D2 is reading the old parity still stored on its disk. After drive D2 has read the old parity, drive D2 executes the XPWRITE command and updates the parity 28.\nDrives D2 and D4 are executing the commands in the order in which the commands are received. In so doing, the old parity and the updated data are read for the regeneration of a data block producing the incorrect data.\nFIGS. 2 through 4 show further examples where competing commands from the same or different hosts for blocks in the same parity stripe can result in inconsistencies.\nFIG. 2 shows REGENERATE and XDWRITE commands as described in the previous example. A host issues a REGENERATE command to a surrogate drive D1 to \"read\" a data block from a failed drive D5. Drive D1 issues READ commands to drives D2, D3, D4 and D5 to read the other blocks of the parity stripe for the requested block 30. At the same time, the same or different host is updating a data block from the same parity stripe stored on drive D2 32. Drive D2 updates the data block and issues an XPWRITE command to drive D4 34. Due to the uncertainty of when the commands will be executed, the READ operation for the reconstruction may return old data or updated data from drive D2 and old parity or updated parity from drive D4.\nReferring to FIG. 3, a REBUILD command is issued to drive D1 which is a repaired drive. Drive D1 issues READ commands to drives D2, D3, D4 and D5 to regenerate a data block B1 36. At the same time, drive D2 receives an XDWRITE command to update a data block B2 needed for the rebuild command. The corresponding parity for blocks B1 and B2 is stored on drive D4. Drive D2 issues an XPWRITE command to drive D4 to update the parity 38. Inconsistencies can arise for the REBUILD command at drives D2 and D4 since the READ operation for the REBUILD at drive D2 may return old data or updated data while the READ operation for the REBUILD being executed at drive D4 may return old parity or updated parity.\nInconsistencies can also arise for an update write to a failed drive. Referring to FIG. 4, a data block B4 on a failed drive D5 is updated. Data block B4 is being written to a surrogate drive D1. In order to determine the updated corresponding parity, surrogate drive D1 acting as a first initiator issues READ commands to drives D2 and D3 for the other data blocks of the parity stripe 40. Drive D1 then writes the updated parity to drive D4 42. At the same time, drive D2 receives an XDWRITE for a data block in the same parity stripe and acts as a second independent initiator for this operation. The read operation issued from drive D1 to drive D2, may return old data or updated data. The parity write operation on drive D4 may be based on one or both data updates and may occur before or after the XPWRITE issued from drive D2.\nThe current proposed ANSI standard is predicated on the assumption that the host will ensure that no data is involved in an update and a reconstruction operation at the same time by coordinating the issuance of commands and reserving drives. This host action adds complexity to an already complex host program especially if two or more hosts are involved. A host independent solution is needed for a drive controlled array that uses XOR commands.\nAccordingly, there is a need to provide a host independent procedure for reconstructing data using array commands on the drives, where each drive can act as temporary initiator for each separate command. It is necessary that such a system avoid the problem of having new data combined with old parity or old data with new parity for a reconstruction operation. Generally, there is a need to coordinate reconstruction and update commands in a RAID system having two or more initiators implementing the RAID functions, such as a system having multiple independent RAID controllers."}
{"text": "1. Field of the Invention\nThis invention relates generally to drain basin inserts, and more particularly to such drain basin inserts that significantly improve the water quality of the water which passes therethrough.\n2. Description of the Related Art\nThis invention is directed to the providing of a drain basin for use with storm sewers. Phase II of EPA's National Discharge Elimination System (NPDES) requires all but the smallest municipal and industrial storm sewer systems to treat storm water discharge to the “maximum extent practicable”. Consequently, the quality of storm water ultimately discharged into storm sewer systems must now meet higher standards than in recent years past. Two pollutants of major concern are sediments and oils that flow into sewer drains during storms, in addition to metals, such as zinc, and nitrogen and phosphorous-containing compounds.\nSome attempts have already been made to address this problem. They include the devices shown in U.S. Pat. Nos. 5,820,762 and 6,217,757, marketed by Hydro Compliance Management of Brighton, Mich. and Revel Environmental Manufacturing known as REM Filters respectively. Another product designed to address the problem is Aquaguard marketed by Aquashield Inc. of Chattanooga, Tenn.\nSeveral patents of interest include Schilling et al, U.S. Pat. No. 6,086,758, which discloses a filtering device for use in storm drains; Leckner et al, U.S. Pat. No. 5,985,157, which discloses a filter device for removing suspended solids from storm water; Bamer et al, U.S. Pat. No. 5,280,762, which discloses a filter insert for a storm drain; Fleischmann, U.S. Pat. No. 6,217,757, which discloses a storm drain filter with vertical screens; Morris et al, U.S. Pat. No. 6,231,758, which discloses a curb-inlet storm drain system for filtering trash and hydrocarbons; Gannon et al, U.S. Pat. No. 6,485,630, which discloses a separation of hydrocarbons from a hydrocarbon-containing liquid; and Allard, U.S. Pat. No. 6,551,023, which discloses a soft bodied high capacity catch basin filtration system.\nIt is thus apparent that the need exists for a drain basin insert which discharges water that has its quality significantly improved, and that can be easily cleaned of sediment, debris, and hydrocarbons and oil-bound pollutants."}
{"text": "Angiogenesis, the formation of new blood vessels, occurs not only during embryonic development and normal tissue growth and repair, but is also involved in the female reproductive cycle, establishment and maintenance of pregnancy, and repair of wounds and fractures. In addition to angiogenesis that occurs in the normal individual, angiogenic events are involved in a number of pathological processes, notably tumor growth and metastasis, and other conditions in which blood vessel proliferation is increased, such as diabetic retinopathy, psoriasis and arthropathies. In addition, angiogenesis is important in the transition of a tumor from hyperplastic to neoplastic growth. Consequently, inhibition of angiogenesis has become an active cancer therapy research field.\nTumor-induced angiogenesis is thought to depend on the production of pro-angiogenic growth factors by the tumor cells, which overcome other forces that tend to keep existing vessels quiescent and stable. The best characterized of these pro-angiogenic agents or growth factors is vascular endothelial growth factor (VEGF) (Cohen et al., FASEB J., 13: 9-22 (1999)). VEGF is produced naturally by a variety of cell types in response to hypoxia and some other stimuli. Many tumors also produce large amounts of VEGF, and/or induce nearby stromal cells to make VEGF (Fukumura et al., Cell, 94: 715-725 (1998)). VEGF, also referred to as VEGF-A, is synthesized as five different splice isoforms of 121, 145, 165, 189, and 206 amino acids. VEGF121 and VEGF165 are the main forms produced, particularly in tumors (see Cohen et al. 1999, supra). VEGF121 lacks a basic domain encoded by exons 6 and 7 of the VEGF gene and does not bind to heparin or extracellular matrix, unlike VEGF165. Each of the references cited in this paragraph is incorporated by reference in its entirety.\nVEGF family members act primarily by binding to receptor tyrosine kinases. In general, receptor tyrosine kinases are glycoproteins having an extracellular domain capable of binding one or more specific growth factors, a transmembrane domain (usually an alpha helix), a juxtamembrane domain (where the receptor may be regulated, e.g., by phosphorylation), a tyrosine kinase domain (the catalytic component of the receptor), and a carboxy-terminal tail, which in many receptors is involved in recognition and binding of the substrates for the tyrosine kinase. There are three endothelial cell-specific receptor tyrosine kinases known to bind VEGF: VEGFR-1 (Flt-1), VEGFR-2 (KDR or Flk-1), and VEGFR-3 (Flt4). Flt-1 and KDR (also known as VEGFR-2 or Flk-1, which are used interchangeably herein) have been identified as the primary high affinity VEGF receptors. While Flt-1 has higher affinity for VEGF, KDR displays more abundant endothelial cell expression (Bikfalvi et al., J. Cell. Physiol., 149: 50-59 (1991)). Moreover, KDR is thought to dominate the angiogenic response and is therefore of greater therapeutic and diagnostic interest (see Cohen et al. 1999, supra). Expression of KDR is highly upregulated in angiogenic vessels, especially in tumors that induce a strong angiogenic response (Veikkola et al., Cancer Res., 60: 203-212 (2000)). The critical role of KDR in angiogenesis is highlighted by the complete lack of vascular development in homozygous KDR knockout mouse embryos (Folkman et al., Cancer Medicine, 5th Edition (B. C. Decker Inc.; Ontario, Canada, 2000) pp. 132-152).\nKDR (kinase domain region) is made up of 1336 amino acids in its mature form. The glycosylated form of KDR migrates on an SDS-PAGE gel with an apparent molecular weight of about 205 kDa. KDR contains seven immunoglobulin-like domains in its extracellular domain, of which the first three are the most important in VEGF binding (Cohen et al. 1999, supra). VEGF itself is a homodimer capable of binding to two KDR molecules simultaneously. The result is that two KDR molecules become dimerized upon binding and autophosphorylate, becoming much more active. The increased kinase activity in turn initiates a signaling pathway that mediates the KDR-specific biological effects of VEGF.\nThus, not only is the VEGF binding activity of KDR in vivo critical to angiogenesis, but the ability to detect KDR upregulation on endothelial cells or to detect VEGF/KDR binding complexes would be extremely beneficial in detecting or monitoring angiogenesis.\nIt is well known that gas filled ultrasound contrast agents are exceptionally efficient ultrasound reflectors for echography. Such ultrasound contrast agents include, for example, gas-filled microvesicles such as gas-filled microbubbles and gas filled microballoons. Gas filled microbubbles are particularly preferred ultrasound contrast agents. (In this disclosure the term of “microbubble” specifically designates a gaseous bubble surrounded or stabilized by phospholipids). For instance injecting into the bloodstream of living bodies suspensions of air- or gas-filled microbubbles in a carrier liquid will strongly reinforce ultrasonic echography imaging, thus aiding in the visualization of internal anatomical structures. Imaging of vessels and internal organs can strongly help in medical diagnosis, for instance for the detection of neoplastic, cardiovascular and other diseases.\nFor both diagnostic and therapeutic purposes it would be particularly beneficial to incorporate into gas filled ultrasound contrast agents, targeting vector compositions which exhibit high binding affinity for a desired target (such as, for example, KDR or the VEGF/KDR complex). For example, targeting vector-phospholipid conjugates and particularly targeting peptide-phospholipid conjugates may be used to prepare targeted, gas filled ultrasound contrast agents. In addition, it would be particularly beneficial to have methods for large scale production of highly purified forms of such targeting vector-phospholipid conjugates. Such compositions and methods would allow for production of compositions for use in diagnostic or therapeutic applications such as, for example, precise targeting of reporter moieties, tumoricidal agents or angiogenesis inhibitors to the target site."}
{"text": "In many applications, the save coupling of connectors is of high importance. For example, in the case of car safety systems, as e.g. airbag systems in passenger cars, the connectors used for the connection of an airbag to its ignition base have to be provided with reliable safety systems. To ensure that the connectors cannot become loose unintentionally, secondary locking members are used to guarantee the save mechanical coupling.\nAn example of a connector with a secondary locking member is described in the WO 97/41623 A1. This document describes a connector which can be mated with a corresponding counter connector being part of an airbag ignition mechanism. In assembled condition, when the connector is mated with the corresponding counter connector, it is fixed to the counter connector by means of flexible latching arms. During mating of the connectors, these arms are deflected and snap back into corresponding latching clearances of the counter connector, when fully mated. For securing the mechanical coupling of the connectors, the WO'623 suggests a secondary locking member which can be inserted into the connector when the two connectors are mated. The secondary locking member comprises locking arms, which upon insertion of the secondary locking member move between faces of the counter connector housing and the latching arms of the connector. Once the locking arms are inserted, they inhibit bending of the latching arms. Thus, the latching arms cannot bend out of the latching clearances and the mechanical coupling of the connectors is secure. The secondary locking member described in WO'623 is further provided with locking means to fix the secondary locking member in its locking position, which is the above described position in which it secures the connector mating. However, in the case described in the WO'623, it is possible to place the secondary locking member in a half-fitting position, in which it may not function properly and in which it is not fixed to the connector, and therefore can be removed unintentionally.\nA further development of a secondary locking member is disclosed in the co-owned patent application DE 100 05 858 A1. This document discloses a connector with a secondary locking device and a safety spring element, which serves to hold the secondary locking member in a pre-locked position, in which the secondary locking member is mounted to the connector housing but does not hinder mating or unmating of the connector with a corresponding counter connector. The safety spring element described in the DE'858 further provides a so-called go/no go function for the secondary locking member: In the pre-locked position of the secondary locking member, two spring arms of the safety spring element are in contact with stop shoulders of the secondary locking member, thereby preventing a movement of the secondary locking member in direction towards its locking position. Upon mounting of the connector to a corresponding counter connector, the counter connector housing comes into contact with a releasing portion of the connector housing, thereby pressing this releasing portion against the safety spring element. By this pressing action, the safety spring element is deformed and the two spring arms are freed from the stop shoulders of the secondary locking member. Upon consequent moving of the secondary locking member towards its locking position, each spring arm is guided by an inclined channel of the secondary locking member, thereby being bent and biased against the edge of this inclined channel. Shortly before the secondary locking member is placed in its locking position, each spring arm passes a turning point of the inclined channel and snaps into a notch of the secondary locking member, thereby fixing the secondary locking member in its locking position. As long as the secondary locking member is not placed in this locking position, the spring arms are biased against the edges of the inclined channel, and due to the resulting restoring spring force they inhibit the secondary locking member to be placed in a half-fitting position unintentionally: i.e. the spring force causes the secondary locking member to move back into its pre-locked position. The safety spring element thereby provides a go/no go mechanism for the secondary locking member, which assures that the secondary locking member is either placed correctly in its locking position or is moved back into its pre-locked position, but prevents any intermediate position.\nIt is an object of the present invention to improve the state of the art by providing a connector assembly with an improved secondary locking mechanism. It is in particular an object of the present invention to provide a connector assembly with a secondary locking member which is provided with an improved and more reliable go/no go mechanism, which does not become non-serviceable, e.g. due to fatigue damage of spring components. It is an even further object of the present invention to accomplish said tasks by providing an inexpensive and non-complex secondary locking member, which preferably can be produced as an injection molded part. These and other objects which become apparent upon reading the following description are solved by a connector assembly according to claim 1."}
{"text": "The present invention relates to a plug part for establishing an optical and/or electrical plug-in connection with a mating plug part by plugging together the plug part with the mating plug part in at least one plugging-together direction, the plug part having a plug part housing and at least one optical and/or electrical contact and an elastic opening prestressing element, acting in a direction parallel to the plugging-together direction, and an opening element for opening a flap on the mating plug part, the opening element being mounted displaceably in relation to the plug part housing in directions parallel to the plugging-together direction in or on the plug part housing, and being prestressed in one of these directions parallel to the plugging-together direction by the opening prestressing element.\nMany configurations of plug parts for establishing optical and/or electrical plug-in connections are known. Not only optical and/or electrical signals but also power can be transmitted by way of the optical and/or electrical plug-in connections. In the case of more sensitive contacts, it is already known in the prior art to protect them from external influences in the unplugged state of the plug part and the mating plug part by a flap. As an example, reference is made here to EP 2 146 232 A1. In this document, both the plug part and the mating plug part each have a flap protecting the contacts, which is only opened during the plugging together of the plug part and the mating plug part."}
{"text": "Adhesive compostions based on alpha-cyanoacrylate esters belong to a class of adhesives known as reactive liquid adhesives. These cyanoacrylate adhesives are single-part, low viscosity adhesives which polymerize rapidly at room temperature without the use of an added catalyst when pressed between two substrates. Application of the cyanoacrylate adhesives merely involves spreading a small amount of the adhesive between two substrates, pressing the substrates together, and allowing the resultant bond to cure. The adhesive develops sufficient strength after a short period of time to hold the substrates together until the adhesive completely polymerizes and builds up to its maximum bonding strength.\nInitiation of polymerization, or cure, is generally believed to proceed through an anionic mechanism. The cyanoacrylate adhesives have such a great tendency to polymerize that water is a sufficiently active initiator. When the adhesive is applied to a substrate and exposed to atmospheric and surface moisture, polymerization usually begins within a relatively short period of time, generally less than one minute, and on many surfaces within a matter of a few seconds.\nFor many years a need has existed for improved bonding of polymeric substrates to other materials, e.g., metals, wood and ceramics, or to themselves with cyanoacrylate adhesives. Untreated polymeric substrates, particularly polyolefin and polyfluorocarbon substrates, when bonded with cyanoacrylate adhesives, generally exhibit less shear strength than is desired for many end uses.\nModification of polymeric surfaces has been used as a means of increasing the adhesion of various adhesives to polymeric surfaces. Surface modification techniques which have been examined are corona discharge treatment, flame treatment with, for example, helium gas plasma or oxygen gas plasma, and chemical treatment with, for example, chromic acid, potassium permanganate, or peroxydisulphate solutions. Such treatments are described by D. Briggs in \"Surface Treatments for Polyolefins,\" Surface Analysis and Pretreatment of Plastics and Metals, D. M. Brewis, Ed., MacMillan, N.Y., 1982, pp. 199-226. These treatments and other surface modification procedures have a common shortcoming in poor durability. Light rubbing of the surface causes a decrease in the effect, the altered surface being easily abraded.\nIrradiation of polyolefin substrates, such as with an electron beam, to improve the adhesion of various coatings is also known and has been disclose, for example, in U.S. Pat. Nos. 4,041,192 (Heger et al.), 4,148,839 (Fydelor), 3,252,880 (Magat et al.) and 4,179,401 (Garnett et al.).\nApplication of a precoat, or primer, on surfaces of various materials has also been used to improve the adhesion of cyanoacrylate adhesives to various substrates. Primers which have been used include alkyl monohydric alcohols (U.S. Pat. No. 2,768,109, Coover, Jr.), monomeric epoxides (U.S. Pat. No. 3,259,534, Wicker, Jr. et al.), organic amines such as secondary amines, tertiary amines, N-substituted alkanolamines, acrylated N-substituted alkanolamines, diamines and certain heterocyclic amines (U.S. Pat. No. 3,260,637, Von Bramer), organometallic compounds (European Patent Publication No. 0 129 069, Toagosei), a solid mixture of alkaline or basic accelerator, such as sodium or potassium hydroxide or sodium borosilicate spheres of micro-size, combined in a chalkable calcareous binder (U.S. Pat. No. 4,215,173, Hubbard), alkyl 2-cyanopenta-2,4-dienoate (U.S. Pat. No. 4,425,471, Millet), and tannins, such as digallic acid, tannic acid, or other hydrolyzable or condensed tannin (U.S. Pat. No. 4,511,686, Millet). However, many of these priming materials are easily removed from the substrate to which they are applied by evaporation or abrasion during handling of the primed substrate.\nImproved adhesion of cyanoacrylate adhesives to various substrates through the use of adhesive additives has also been described. Various additives which have been suggested include plasticizers, such as alkyl esters of aliphatic monocarboxylic acids, alkyl esters of aliphatic dicarboxylic acids, alkyl phosphates, triaromatic phosphates, polyfunctional aliphatic esters, and aliphatic and aromatic phosphonates (U.S. Pat. No. 2,784,127, Joyner et al.), itaconic acid anhydride (U.S. Pat. No. 3,948,794, Eberhard), acetic acid (U.S. Pat. No. 4,125,494, Schoenberg et al.), compounds having the formula ##STR1## wherein R' is hydrogen or an alkyl, aryl, or cycloalkyl group having 1 to 10 carbon atoms, such as gallic acid monohydrate, methyl gallate, propyl gallate and hexyl gallate (U.S. Pat. No. 4,139,693, Schoenberg), a combination of polyethylene glycols having a degree of polymerization of at least 3 and non-ionic surface active agents having a poly(ethyleneoxy) moiety therein with the poly(ethyleneoxy) moiety having a degree of polymerization of at least 3 (U.S. Pat. No. 4,170,585, Motegi et al.), trichlortrifluoroethane (U.S. Pat. No. 4,200,549, Okamura et al.), esters compounds such as acrylates, methacrylates, and crotonates of glycols, glycol monoethers, and monoesters (U.S. Pat. No. 4,307,216, Shiraishi et al.), a combination of (1) at least one aliphatic polyol and/or polyether and derivatives thereof and (2) at least one aromatic polyol and/or carboxylic acid and derivatives thereof (U.S. Pat. No. 4,377,490, Shiraishi et al.), and phthalic anhydride (U.S. Pat. No. 4,450,265, Harris).\nPolyolefin substrates have been provided with a treated surface for improved adhesion of pressure-sensitive adhesives. U.S. Pat. No. 3,628,987 (Nakata et al.) discloses a pressure-sensitive adhesive film wherein the film surface to which the adhesive is adhered has graft-polymerized thereto a vinyl monomer or diene monomer, the adhesive having a solubility parameter near that of polymers of the vinyl or diene monomer. U.S. Pat. No. 4,563,388 (Bonk et al.) discloses a polyolefin substrate having graft-polymerized thereto at least one monomer selected from the group consisting of acrylic acid, methacrylic acid and ester thereof; acrylamide; methacrylamide; sterically non-hindered tertiary alkyl acrylamides and methacrylamides having three or less carbon atoms in the alkyl group; and N-vinyl pyrrolidone, and firmly adherently bonded to the graft-polymerized monomer, an acrylic-type, normally tacky and pressure-sensitive adhesive.\nA series of articles, \"Surface Modification of Polyethylene by Radiation-Induced Grafting for Adhesive Bonding. I. Relationship Between Adhesive Bond Strength and Surface Composition,\" (S. Yamakawa, J. Appl. Polym. Sci., 20, 3057-3072 (1976); \"II. Relationship Between Adhesive Bond Strength and Surface Structure,\" (S. Yamakawa et al., Macromolecules, 9, 754-758, 1976); \"III. Oxidative Degradation and Stabilization of Grafted Layer,\" (S. Yamakawa et al., J. Appl. Polym. Sci., 22, 2459-2470, 1978); \"IV. Improvement in Wet Peel Strength,\" (S. Yamakawa et al., J. Appl. Polym. Sci., 25, 25-39, 1980), and \"V. Comparison with Other Surface Treatments,\" (S. Yamakawa et al., J. Appl. Polym. Sci., 25, 40-49, 1980), disclose grafting of methyl acrylate (followed by saponification), vinyl acetate, acrylic acid, acrylamide, and methylolacrylamide to polyethylene by vapor-phase mutual grafting or liquid-phase preirradiation at thickness of grafted monomer of more than 10 micrometers to improve adhesion of epoxy adhesives."}
{"text": "In digital signal analysis, a signal is typically reconstructed from discrete measurements. Typical approaches to sampling signals or images follow Shannon's theorem wherein a sufficient sampling rate is at least twice the bandwidth (i.e., a lowpass or bandpass bandwidth depending on the signal type) present in the signal. This minimal sampling rate is known as the Nyquist rate or frequency.\nThis principle underlies nearly all signal acquisition protocols used in consumer audio and video electronics, medical imaging devices, radio receivers, etc. In reference to data conversions, for example, a standard analog-to-digital converter (ADC) implements the quantized Shannon representation in that the signal is uniformly sampled at or above the Nyquist rate.\nEven though traditional approaches to data acquisition depend on a Nyquist rate, in many applications the Nyquist rate far exceeds the necessary sample rate needed (or the number of samples over the time record) to accurately reconstruct the signal. For efficient storage or communications of the signal information content, data compression is an important processing step prior to storage or transmission.\nHowever, there are penalties associated with implementing data compression. One such penalty is that the original information bearing signal will likely be sampled above the Nyquist rate, and thus produce many samples with redundant information content about the signal. This process can at times require an ADC with a very high sample rate. High sample rate converters tend to require higher power and are often limited in resolution with respect to their lower sample rate counterparts. This induces a system penalty of low resolution and high power. In some cases it is difficult to even purchase or construct an ADC with the required sample rate given the Nyquist viewpoint of the signal.\nAnother penalty is induced by the data compression step, as performed with a Karhunen-Loève transform, for example. The data compression step attempts to remove the information redundancy in the samples, induced by the Nyquist view of the conversion, so that a minimal set of samples result to sufficiently represent the information content of the original signal. Implementing compression prior to storage or transmission requires some type of computing resource, which in turn increases the power drawn by the sensor system. The larger the oversampling factor in a Nyquist paradigm, the more severe these penalties become.\nRecent developments have shown that compressive sampling or compressive sensing can provide sub-Nyquist rate sampling for communications systems. One such approach is disclosed in U.S. Patent No. 2011/0090394 to Tian et al. A disclosed method of signal processing includes receiving at a processor a data packet comprising compressively sensed or measured data of a signal, with the compressively measured data comprising wavelet transform coefficients. The received signal is a discrete signal, which in turn, requires transform coding before being processed by the processor. The processor reconstructs the signal using a clustering property of the wavelet transform coefficients. A disadvantage of this approach is that transform coding of the discrete signal before being compressively sampled requires additional processing, which in turn consumes power. Additional power consumption may be undesirable, particularly for battery-powered systems."}
{"text": "1. Field of the Invention\nThis invention relates to a ferrite magnet suitable for use in electronic copying apparatus, audio equipment, rotary machines, etc., and a method of producing such ferrite magnet by means of extrusion compacting.\n2. Description of the Prior Art\nHeretofore, a ferrite magnet used in speakers, small DC motors, magnet type rotary machines, etc., has usually been produced by compacting, and when an anisotropic ferrite magnet of high magnetic properties is produced, magnetic field compacting in which compacting is carried out in the magnetic field has been in use. For compacting a cylindrical or columnar ferrite magnet of elongated type (having a length at least twice as large as the diameter) suitable for use as a magnet roll used in a developing device of an electronic copying apparatus, for example, what is generally referred to as a rubber press method is used which consists in compacting a ferrite magnet by applying pressure thereto at a right angle to the axis. Such magnet roll is described in specifications of U.S. Pat. Nos. 4,166,263, 4,168,481, 4,266,328, 4,167,718 and 4,169,998 (granted to Yamashita et al.). However, it is only an isotropic ferrite magnet that is produced by the rubber roll method.\nIn many applications, an ordinary magnet roll is required to have the magnetic properties of Br.gtoreq.3300 G and of the magnetic flux density Bo.gtoreq.700 G on the sleeve. In addition, it is also required that a compact size be obtained (having an outer diameter of below 24 mm) in a magnet roll. In the aforesaid rubber press method, however, difficulties would be experienced in obtaining improved uniformity of the magnet in a major axis direction because of the fact that compacting is carried out through a rubber mold. Another disadvantage in the rubber press method is that when sintered bodies of about 25 mm in diameter are ground to a size of about 24 mm in outer diameter, the Br is in the range 2000 and 2100 G at most because the ferrite magnet obtained in isotropic.\nIn one method known in the art for producing an anisotropic ferrite magnet, anisotropy is imparted beforehand prior to compacting, as described in Japanese Patent Publication No. 47043/78, for example. Another known method is an extrusion compacting which is used for producing a cylindrical rubber magnet or a plastic magnet used in a small motor, for example. It is disclosed in Japanese Patent Application Laid-Open No. 75536/74 that an elongated cylindrical ferrite magnet for use as a magnet roll is produced by extrusion compacting.\nThe extrusion compacting methods referred to hereinabove have high forming efficiency, so that productivity can be greatly increased. Besides they are capable of imparting anisotropy to a certain degree to the magnet when it is produced, thereby enabling an improvement in magnetic properties to be expected. However, the ferrite magnet produced by the extrusion compacting methods of the prior art is unable to meet the properties as an anisotropic ferrite magnet. To obviate this disadvantage of the extrusion compacting methods of the prior art, it has been usual practice to carry out extrusion compacting in a magnetic field and it is proposed, as shown in Japanese Patent Publication No. 34192/72, to use a special nozzle in addition to the application of a magnetic field. Thus, compacting of a ferrite magnet in a magnetic field is generally practised to obtain an anisotropic ferrite magnet by extrusion compacting. One encounters the problems, when this method is used, that a static magnetic field of a current of a high value is required and hence the apparatus becomes large in size, and in addition the magnetic field must have a considerably high value of about 5000 G."}
{"text": "(1) Field of the Invention\nThe present invention relates to an image processing method, an image processor, an integrated circuit, and a recording medium for generating an output image of a resolution higher than a resolution of an input image, using the input image.\n(2) Description of the Related Art\nWhen a resolution of an image to be displayed on a high resolution display is lacking, it is necessary to enlarge the image so as to match the resolution of the image to a resolution of the high resolution display.\nAlthough methods of increasing the resolution of an image have been conventionally suggested, there is a limit on processing capacity for the practical use. Thus, the combination of simple image enlargement by interpolation and image enhancement processing has supported the methods. Thus, problems of image degradation have occurred, such as blurring and noticeable jaggies in an edge.\nRecent technological advances in performance of hardware has enabled allocation of larger processing capacity to the image enlargement processing. Here, attention is currently focused on the super-resolution technique capable of converting an image to an image with high quality and high resolution through complex processing.\nIn the super-resolution technique, a method using a training database storing data learned from examples of correspondence between high-resolution images and low-resolution images is called a training-based super-resolution.\nFIG. 21 is a block diagram illustrating a configuration of an image processor 600 that performs the same processing as in the conventional training-based super-resolution disclosed in Patent Reference 1. Hereinafter, the technique disclosed in Patent Reference 1 (Japanese Unexamined Patent Application Publication No. 2003-018398) is also referred to as a conventional technique A.\nAs illustrated in FIG. 21, the image processor 600 includes a training database 610, an image enlargement unit 605, a search vector generating unit 620, a training data search unit 630, and an addition unit 680.\nNext, operations of the image processor 600 will be described.\nFIG. 22 is a flowchart of the procedure performed by the image processor 600 in FIG. 21.\nThe image processor 600 converts an input image 601 of a low resolution into an output image 602 of a high resolution, using the training database 610 that stores index vectors and high-frequency component data items 631 that are respectively associated with the index vectors. Each of the index vectors is an image of a low-frequency component.\nThe training database 610 stores the index vectors and the high-frequency component data items 631 that are respectively associated with the index vectors. Each of the index vectors and the high-frequency component data items 631 is a block-shaped image.\nAt Step S4001, the image enlargement unit 605 generates a temporary enlarged image 606 by enlarging the input image 601.\nThe following processes are performed on each target block included in the input image 601.\nAt Step S4002, the search vector generating unit 620 generates a search vector 621 by concatenating a medium-frequency component data item of a target block extracted from the temporary enlarged image 606, with a high-frequency component data item of an adjacent block which has been processed. In other words, the search vector 621 is composed of the medium-frequency component data item and the high-frequency component data item in an overlap region with the adjacent block. The search vector 621 is a block-shaped image.\nAt Step S4003, the training data search unit 630 selects one of the high-frequency component data items 631 that corresponds to an index vector having the highest similarity to the search vector 621, from among the index vectors stored in the training database 610.\nAt Step S4004, the addition unit 680 adds the selected high-frequency component data item 631 to the temporary enlarged image 606. Thereby, an output image 602 of a high resolution is generated.\nIn the conventional training-based scheme, one of the high-frequency component data items 631 that corresponds to the index vector having the highest similarity to the search vector 621 is searched. Since only one of the high-frequency component data items 631 is selected from the database in this method, when the high-frequency component data item 631 selected from the training database 610 indicates noise, the noise immunity is low.\nThus, since a bumpy edge and artifacts occur in the output image 602 of the high resolution, the image quality is degraded. Furthermore, since the search vector 621 includes the overlap region included in a high-frequency component and in the adjacent region, there is a possibility of error propagation."}
{"text": "1. Field of the Invention\nThe present invention relates to a lead pin for a package substrate, and more particularly, to a lead pin for a package substrate, which connects the package substrate, on which an integrated circuit (IC) is mounted, to a main board.\n2. Description of the Related Art\nAs the electronics industry has developed, various types of semiconductor packages have been manufactured. As interconnections for semiconductor packages have become highly dense, a Pin Grid Array (PGA) semiconductor package substrate on which a plurality of T-type lead pins are installed is widely used as a substrate which connects a package substrate, on which an IC is mounted, to a main board.\nA typical package substrate generally uses an insertion-type pin, which is inserted into a through-hole, and a T-type lead pin, which is attached to a package substrate by soldering. Compared with the insertion-type pin, the T-type lead pin is commonly used because there is less limitation on the circuit configuration of the package substrate.\nSince the use of Pb has recently been restricted due to the consideration of environmental impact of soldering, a Pb-free solder such as Sn—Ag—Cu or Sn—Sb is used. Thus, a solder melting temperature is high.\nAs the solder melting temperature increases, a solder which supports a lead pin may be melted by reflow heat during a reflow process for mounting an IC chip on a package substrate, causing the lead pin to be inclined.\nFurthermore, a solder often flows over the top surface of the lead pin and then flows toward the coupling pin during the reflow process. Therefore, there is a need for technologies which can solve these problems."}
{"text": "The invention relates to improvements in dryers, and more particularly to a rotary dryer of the type which may be used for drying bulk material.\nMore particularly, the invention relates to dryers of the type which may be used in papermaking plants for drying hog fuel or bagasse with waste heat and direct fired dryers used for difficult to dry materials such as pulp mill sludge. The interest in drying systems for drying materials to be used as fuel results from the fact that dry fuel increases boiler efficiency and fuel BTU value. In the case of sludge, drying is beneficial in preparing the wet material as fuel or to substantially reduce its bulk, thereby reducing or eliminating costly landfill requirements. If drying of fuel is done in the boiler, much of the energy of combustion which could be used to generate steam is consumed by the drying process of incoming fuel. Because boilers are inefficient dryers, the effect on performance is dramatic in that efficiency declines substantially where the fuel must be dried within the boiler. In the case of sludge, drying in the boiler is extremely disruptive to the boiler operation.\nIn the boiler, wet fuels require large amounts of excess air to sustain combustion. The excess air, combined with the water vapor generated during drying reduces boiler efficiency and furnace temperature. The colder furnace produces less steam and more particulate emissions as unburned fuel passes up the stack.\nDrying fuel results in more BTU's per pound of fuel as well as an increased boiler efficiency. Drying also allows the effective control of fuel moisture content as the moisture varies by season and with the composition of the raw material. The consistency of the fuel also contributes to more efficient and predictable boiler performance.\nFuel dryers have been used heretofore which are heated by waste heat and by direct firing, but the efficient and effective transfer of heat from the hot gases to the wet material varies greatly depending on the material properties and on the dryer loading.\nIt is accordingly an object of the present invention to provide an improved dryer such as may be used for drying hog fuel, sludge or bagasse with waste heat and by furnace direct firing wherein a unique and improved transfer of heat from the hot gases to the wet material is accomplished.\nA further object of the invention is to provide a rotary dryer wherein the contents are tumbled or moved in a unique manner so as to increase the efficiency, capacity and effectiveness of heat transfer to the material even at low feed rates.\nA further object of the invention is to provide an improved flight system within a rotary dryer wherein panels are provided arranged so that improved heat transfer results through controlling gas swirl and mixing, and wherein flights are provided having metal surfaces which effectively increase heat transfer through conduction.\nAnother object of the present invention is to provide an improved rotary drum center flight design which has a significant return effect on material passing through the dryer to significantly counteract the gas flow forces increasing retention time and drying efficiency on difficult to dry materials such as sludge."}
{"text": "The optical fiber communication, which was mainly used for business, has also been widely used for home use. As a result, an optical communication device of high performance has been required. As an optical communication device for use in various optical communication systems, such as an optical fiber communication system for home use and a local area network (LAN) system, there is a silicon-based optical communication device which functions at optical fiber communication wavelengths of 1330 nm and 1500 nm. The silicon-based optical communication device is a very promising device in which an optical functional element and an electronic circuit can be integrated on a silicon platform by using a CMOS (Complementary Metal Oxide Semiconductor) technique.\nAs silicon-based optical communication devices, passive devices, such as a waveguide, an optical coupler, and a wavelength filter, have been widely investigated. Further, active devices, such as a silicon-based optical modulator and a silicon-based optical switch, are listed as important elements of means for operating an optical signal in the above-described optical communication system, and have been attracting much attention. The optical modulator and the optical switch, in which the refractive index is changed by using the thermo-optic effect of silicon, have a low optical modulation speed, and hence can be used only for a device having an optical modulation frequency of 1 Mb/s or lower. An optical modulator using the electro-optic effect is required for a device having an optical modulation frequency higher than 1 Mb/s.\nPure silicon does not exhibit any change in its refractive index due to the Pockels effect and exhibits a very small change in its refractive index due to the Franz-Keldysh effect and the Kerr effect. For this reason, many of the optical modulators that is proposed at present and that use the electro-optic effect utilize the carrier plasma effect. That is, the phase and intensity of light are changed by changing the real part and imaginary part of the refractive index by changing the density of free carriers in a silicon layer.\nThe density of free carriers in the optical modulator can be changed by the injection, accumulation, depletion or inversion of the free carriers. Many of such optical modulators, which have been investigated to date, have a low optical modulation efficiency, and require, for optical phase modulation, a length of 1 mm or more and an injection current density of 1 kA/cm3 or more. In order to realize miniaturization, high integration, and low power consumption in the optical modulator, a device structure having high optical modulation efficiency is required. When a device structure having high optical modulation efficiency is realized, the length required for the optical phase modulation can be reduced. Further, in the case where an optical communication device has a large size, it is also conceivable that the optical communication device is liable to be affected by the temperature on the silicon substrate, and thereby the electro-optic effect to be originally obtained is cancelled by a change in the refractive index of the silicon layer resulting from the thermo-optic effect.\nFIG. 1 shows an example of a related art of a silicon-based electro-optical optical modulator using a rib waveguide formed on an SOI (Silicon on Insulator) substrate. Embedded oxide layer 32 and intrinsic semiconductor 31 including a rib-shaped portion are laminated in order on substrate 33. P+-doped semiconductor 34 and n+-doped semiconductor 35 are respectively formed on both sides of the rib-shaped portion of the intrinsic semiconductor 31 so as to be separated from each other by a distance. P+-doped semiconductor 34 and n+-doped semiconductor 35 are formed by performing a high concentration doping process to respective parts of intrinsic semiconductor 31. The optical modulator shown in FIG. 1 is configured as a PIN (P-intrinsic-N) diode. When a forward bias voltage or a reverse bias voltage is applied to the PIN diode, the density of free carriers in intrinsic semiconductor 31 is changed, so that the refractive index is changed by the carrier plasma effect. In this example, electrode contact layer 36 is arranged on one side of the rib-shaped portion of intrinsic semiconductor 31, and p+-doped semiconductor 34 described above is formed at a position facing electrode contact layer 36. Similarly, electrode contact layer 36 is arranged also on the other side of the rib-shaped portion of intrinsic semiconductor 31, and n+-doped semiconductor 35 is formed at a position facing electrode contact layer 36. Further, the waveguide including the rib-shaped portion is covered by oxide clad 37. In the above-described PIN diode structure, the doping process can be performed at a high concentration so that the density of carriers in semiconductors 34 and 35 becomes about 1020/cm3.\nWhen an optical modulation operation is performed, forward bias voltage is applied to the PIN diode from a power source connected to electrode contact layers 36, so that free carriers are injected into the waveguide. At this time, free carriers are increased, and thereby the refractive index of intrinsic semiconductor 31 is changed, so that the light propagated through the waveguide is phase-modulated. However, the speed of the optical modulation operation is restricted by the lifetime of free carriers in the rib-shaped portion of intrinsic semiconductor 31, and by the carrier diffusion at the time when the application of the forward bias voltage is stopped. The electro-optical optical modulator having the related art PIN diode structure as described above usually has an operating speed in the range from 10 to 50 Mb/s at the time when the forward bias voltage is applied. On the other hand, by introducing impurities into intrinsic semiconductor 31 in order to reduce the lifetime of carriers, the switching speed of the electro-optical optical modulator can be increased. However, there is a problem that the introduced impurities lower the optical modulation efficiency. Further, the major factor that influences the operation speed is the RC time constant. The electrostatic capacitance at the time when the forward bias voltage is applied is significantly increased by the reduction in the carrier depletion layer of the PN (Positive-Negative) junction portion. Theoretically, the high-speed operation of the PN junction portion can be realized by applying a reverse bias voltage. However, a comparatively high drive voltage or a large element size is required in this case.\nAs another example of the related art, Patent literature 1 discloses a silicon-based electro-optical optical modulator having a capacitor structure in which embedded oxide layer 32 and a first conductivity type main body region are laminated in order on substrate 33, in which a second conductivity type gate region is further laminated so as to partially overlap the main body region, and in which thin dielectric layer 41 is then formed at the lamination interface between the main body region and the gate region. Note that hereinafter, “thin” is intended to mean a thickness of submicron order (less than 1 μm).\nFIG. 2 shows a silicon-based electro-optical optical modulator having an SIS (silicon-insulator-silicon) structure according to the related art. The electro-optical optical modulator is formed on an SOI substrate configured by substrate 33, embedded oxide layer 32, and a main body region. The main body region is configured by p-doped semiconductor 38 formed by performing a doping process on the silicon layer of the SOI substrate, p+-doped semiconductor 34 formed by performing a high concentration doping process on the silicon layer of the SOI substrate, and electrode contact layer 36. The gate region is configured by n-doped semiconductor 39 formed by performing a doping process on a thin silicon layer laminated on the SOI substrate, n+-doped semiconductor 35 formed by performing a high concentration doping process on the thin silicon layer, and electrode contact layer 36. Oxide clad 37 is provided in the gap formed by embedded oxide layer 32, the main body region and the gate region, and is also provided on the main body region and the gate region.\nThe region subjected to the doping process is configured such that the change in carrier density is controlled by an external signal voltage. Further, when a voltage is applied between electrode contact layers 36, free carriers are accumulated, depleted or inverted on both sides of dielectric layer 41. Thereby, the optical phase modulation is performed. For this reason, it is preferred that the region of optical signal electric field be made to coincide with the region in which the carrier density is dynamically controlled from the outside."}
{"text": "Implantable medical devices (IMDs) detect and deliver therapy through a lead. Typically, a stimulation electrode at the distal end of a lead is positioned near or in tissue so that electrical stimuli may be delivered. To reduce or prevent inflammation of the tissue in response to the stimulation electrode, the distal end of the lead includes a monolithic controlled release device that releases an anti-inflammatory agent such as a steroid. It is desirable to develop devices that are able to chronically release anti-inflammatory agents in the vicinity of a stimulation electrode."}
{"text": "1. Technical Field\nThe present invention relates to an array substrate for mounting a chip and a method of the same, and more particularly, to an array substrate for mounting a chip in which an optical device chip is mounted.\n2. Description of the Related Art\nAs well known, a liquid crystal display (LCD) device is widely used in a flat panel display unit such as a television, or a computer monitor and the like, and a back light unit (BLU) is an illumination part of emitting light in the back of the LCD device.\nAccordingly, the BLU includes a light guide panel for guiding the light incident from the side surface to the front, and a light emitting diode (LED) array is recently used as a side line source of light in the light guide panel. However, according to a conventional optical device array, when manufacturing a module by combining each of LED packages in an array form, the LED packages have to be classified and packaged in order to minimize color deviation of the LED packages. However, there is a difficulty in a process when classifying and packaging the LED packages divided into each LED package."}
{"text": "Under the Adaptive Multi Rate (AMR) speech codec standard, the codec mode applied to the voice data exchanged between a mobile communication device and a communication network is dynamically selected by the communication network based on changing radio channel conditions and capacity requirements. The codec modes that may be applied during an active voice call are chosen from an Active Codec Set (ACS), which include codec modes defined by the communication network as part of the initial voice call setup. When the communication network determines based on one or more channel quality indicators that radio channel conditions are poor, the communication network may select a lower bit rate codec mode from the ACS in order to accommodate an increase in redundant data added as a result of channel coding. But if the radio channel conditions improve, the communication network may select a higher bit rate codec mode from the ACS, which enhances the perceived voice quality.\nIn general, a conventional mobile communication device will apply the codec mode commanded by the communication network. As such, a codec mode command (CMC) from the communication network may trigger an immediate change (e.g., upgrade or downgrade) in the current uplink (UL) codec mode applied by the mobile communication device. In this respect, codec mode adaptation may in fact compromise the performance of the conventional mobile communication device, particularly in the case of a dual-subscriber identity module (SIM), dual active (DSDA) mobile communication device.\nIn a DSDA mobile communication device, a timeslot collision may occur when both subscriptions simultaneously attempt to perform an activity. To resolve a timeslot collision when both subscriptions attempt to perform an activity (e.g., an UL activity), the transmission of voice data (e.g., speech frames) on one subscription may be inhibited if the voice data transmission has a lower priority transmission than the transmission (e.g., signaling) on the other subscription. Also, the transmission of voice data on one subscription may be inhibited (i.e., blanked) in favor of higher priority receptions on the other subscription. As a result of frequently blanking the transmission of voice data, UL voice data from the DSDA mobile communication device tends to exhibit numerous interruptions. The communication network typically misconstrues these interruptions in the UL voice data from the DSDA mobile communication device to imply fluctuating radio channel conditions. As a result, CMCs from the communication network require the DSDA mobile communication device to change its codec mode at an excessive rate and to unnecessarily downgrade to codec modes with low bit rates, which consequently lowers the Mean Opinion Score (MOS) of the voice call."}
{"text": "Magnetic fusion reactions, such as between deuterium and tritium combined in a hot plasma, have the potential to provide energy with minimal production of long-lived radioisotopes. It is well known that such reactions take place under very high temperatures conditions at which particle momentum is sufficient to overcome the mutual electrostatic repulsion of the reacting nuclei. Physical walls are not feasible for containing the reactions at the required temperatures, because the walls will cool the reactants, thereby quenching the reaction. As such, devices or apparatuses have been pursued that use non-physical means to containing the fusion reactions. These non-physical means of containment have used inertia, magnetism and electrostatic forces as the means for containing the reactions. Because the present invention is in the class of devices that use magnetic forces, the following will be confined to a discussion of this class of device.\nHot plasma is an electrically conductive fluid that will flow with little resistance along magnetic field lines, but diffuse slowly across field lines because it is retarded by the magnetic reaction forces resulting from currents induced in the plasma. For this reason, an objective in designing a magnetic containment system is to suspend the plasma in a field with lines which form closed paths within the plasma confinement volume, and therefore do not provide direct escape paths along field lines. Closed current loops in the conductive plasma form closed magnetic flux loops within the plasma. A simple current loop, for example, forms poloidal flux loops that enclose a toroidal volume which surrounds the current loop. Three magnetic confinement schemes using such internal magnetic currents to generate at least a portion of the confining field are tokamaks, spheromaks and theta pinch systems.\nSpheromaks are formed by two interlinked plasma current loops generated by inductive and/or coaxial gun conductive means, resulting in a near-spherical toroidal confinement volume with field components in both the poloidal and toroidal directions. Tokamaks are toroidal confinement volumes formed by toroidal field components from an external set of poloidal coils and poloidal field components from a toroidal current loop induced by a time-changing external field. Theta pinch systems are toroidal confinement volumes formed by a single plasma current loop induced by a time-changing external field, with field components primarily in the poloidal direction. Reference also should be made to U.S. Pat. No. 4,436,691; Thomas R. Jarboe, \"Review of Spheromak Research\", Plasma Physics and Controlled Fusion, Vol. 36, pp. 945-990 (1994); Kenro Miyamoto, Plasma Physics for Nuclear Fusion (Revised Edition), pp. 530-552, The MIT Press, Cambridge, Mass. (1987); and M. Tuszewski, \"Field Reversed Configurations\", Nuclear Fusion, Vol. 28, No. 11, pp. 2033-2092 (1988).\nThere are a number of known processes for forming currents in the plasma as briefly discussed below. In one process, an inductive transformer action induces transient current loops in a stationary plasma, using time-changing external fields. This process is simple and does not require conductive electrodes to contact with the plasma. However, it is by nature transient, not steady-state. This principle is used in the tokamak and theta pinch approaches\nIn another process, a conductive current transfer forms a plasma current between two electrodes, and then the plasma is moved away from the electrodes. The movement causes the current to detach from the electrodes and reconnect as a closed loop. The process may be repeated rapidly to form a sequence of current loops which merge with and sustain a preexisting current loop. See also Y. Ono, A. Morita, M. Katsurai and M. Yamada, \"Experimental Investigation of Three-Dimensional Magnetic Reconnection by Use of Two Colliding Spheromaks,\" Phys. Fluids B, Vol. 5, pp 3691- to 3701 (1993). This principle is used in the coaxial gun spheromak formation approaches referenced above.\nIn yet another process, an inductive current transfer forms a plasma current in a plasma volume, and then the plasma is moved away from the formation area. The movement causes a process in which magnetic reconnection results in a closed current loop that is no longer inductively linked with the original formation field. As with the conductive current transfer process, this process may be repeated rapidly to form a sequence of current loops which merge with and sustain a preexisting current loop. This principle is used in the inductive and conical theta pinch spheromak formation approaches referenced above.\nParticle beams or traveling waves may be used to differentially move the plasma electrons relative to the positive nuclei and create a steady-state current within the plasma. This method is discussed in the tokamak confinement scheme referenced above. It is also used in the rotamak, confinement scheme, a variation of the spheromak. See also P. M. Bellan, \"Particle Confinement in Realistic 3D Rotamak Equilibria\", Physical Review Letters, Vol. 62, No. 21 pp 2464-2467 (1989).\nRadial diffusion of plasma through the poloidal field lines of a tokamak has been observed to generate a toroidal \"bootstrap\" current which supplements or replaces the inductively generated current (See Kenro Miyamoto, Plasma Physics for Nuclear Fusion (Revised Edition), pp. 225 and 552, The MIT Press, Cambridge, Mass. (1987)). This is a first order electromagnetic dynamo process, and is potentially steady-state, assuming a sustained flow of fresh plasma into the confinement volume. It has the further advantage of being driven by the temperature and pressure differences between the central fusion reaction zone and the outside of the containment volume, and therefore does not consume external electric power.\nMagnetic mirror plasma confinement devices generate open solenoidal magnetic fields with stronger mirror fields at the ends to reflect plasma particles back toward the center of the containment volume. Experiments were carried out on various means of inserting plasma into these devices, including shooting high velocity pulses of plasma axially through the mirrors, e.g., see B. W. Johnson and J. G. Siambis, \"Injection of a Streaming Plasma Into a Mirror Machine,\" Plasma Physics, Vol. 15, pp 369 to 374 (1973). The experimental data include the finding that the peak transient plasma density in the inlet throat of the mirror is about 6 times higher than the peak plasma density achieved at the midplane of the device. The authors interpret this high inlet density as the result of a shock wave, and treat it as a problem to be overcome in achieving efficient plasma injection into the mirror machine.\nThe above described systems, have not progressed to the stage where conditions conducive to fusion reactions have been continuously maintained so that energy can be reliably produced for consumption. As such, there is still need for devices, systems and methods for confining plasmas under conditions conducive to fusion reactions and more particularly a system/method where a steady-state current is formed internal to the plasma."}
{"text": "With developments of mobile devices and information and communication technology, a common authentication provided through a secure element, hereinafter, also referred to as SE, such as a universal subscriber identity module (USIM), and a secure digital (SD) card, are increasingly being applied to an electronic payment.\nHowever, the common SE may manage authentication means such as certificate authentications in a form of a file and thus, may a risk of accidents caused through replication of the certificate authentication or leakage of personal information. To minimize the risk, a one-time password (OTP) authentication may be performed using an existing OTP providing device of a user irrespective of a presence of the SE.\nIn general, an OTP providing apparatus may be provided in a form of, for example, an OTP token, a card type OTP terminal, and a USIM software-based OTP providing application.\nHowever, the OTP providing device may be subject to be exposed to a security attack such as a debugging port and an internal memory attack. Also, the OTP providing device may also be vulnerable to software hacking on a terminal operating system (OS) and OTP applications.\nFurthermore, when a security accident occurs in a certification authority or a service provider that provides electronic payment and mobile banking services, an issue that financial institutions being unable to avoid responsibility for illegal transactions based on leaked information may arise.\nConversely, a physically unclonable function (PUF) may provide an unpredictable digital value. Each PUF may provide a different digital value although the PUF is produced through the exactly same manufacturing process.\nThe PUF may also be referred to as, for example, a physical one-way function (POWF) practically impossible to be duplicated.\nDue to the impossibility of replication, the PUF may be used as an identifier of a device for security and/or authentication. For example, the PUF may be used to provide a unique key for distinguishing a device from other devices.\nThe method of implementing a PUF, Korea Patent No. 10-1139630, hereinafter, referred to as the '630 patent, may be provided in advance. The '630 patent provides a method of stochastically determining whether a via or an inter-layer contact between conductive layers of a semiconductor based on a process variation of the semiconductor."}
{"text": "Field\nEmbodiments disclosed herein generally relate to an apparatus and method for controlling the pressure of a processing chamber. More particularly, embodiments herein relate to a dynamic pressure control system.\nDescription of the Related Art\nSubstrates, such as semiconductor substrates, can be subjected to an epitaxial growth process to form an epitaxial layer on a surface of the substrate. Conventional epitaxial growth processes include flowing a process gas laterally over the surface of the substrate, and thermally decomposing the process gas on the surface of the substrate in order to deposit the epitaxial layer.\nAfter deposition, process gases remaining in the atmospheric pressure epitaxial growth systems are carried out of the processing chamber to facilitate exhaust through an abatement system. Due to this configuration, the pressure in the processing chamber is affected by the surrounding atmospheric pressure and the exhaust side abatement system (i.e., to remove byproduct clogging), thus there is risk for exhaust byproducts to flow back into the processing chamber by pressure oscillations in the processing chamber. Pressure fluctuations, or pressure spikes, in the processing chamber may result in non-uniform distribution of the processing gas. Therefore, maintaining a constant pressure inside the processing chamber is key for uniform deposition results.\nOne option for controlling pressure inside the processing chamber is through the use of a vacuum pump. However, atmospheric epitaxial growth systems configured with a vacuum pump can become unsafe due to faster pyrophoric byproduct build up and risk for explosion and fire in exhaust facilities. Therefore, atmospheric epitaxial growth systems are generally not configured with vacuum pump.\nAnother option for controlling pressure inside the processing chamber includes a pressure control system in the exhaust line to control fluctuations in pressure in atmospheric epitaxial growth systems. Conventional designs can reduce the processing chamber pressure fluctuations to some extent, but conventional designs lack a dynamic range of control.\nThus, there is a need for improved pressure control system processing chambers, and particularly in atmospheric processing chambers."}
{"text": "1. Field of the Invention\nThe present invention relates in general to the field of information handling system power management, and more particularly to a system and method for information handling system adaptive variable bus idle timer.\n2. Description of the Related Art\nAs the value and use of information continues to increase, individuals and businesses seek additional ways to process and store information. One option available to users is information handling systems. An information handling system generally processes, compiles, stores, and/or communicates information or data for business, personal, or other purposes thereby allowing users to take advantage of the value of the information. Because technology and information handling needs and requirements vary between different users or applications, information handling systems may also vary regarding what information is handled, how the information is handled, how much information is processed, stored, or communicated, and how quickly and efficiently the information may be processed, stored, or communicated. The variations in information handling systems allow for information handling systems to be general or configured for a specific user or specific use such as financial transaction processing, airline reservations, enterprise data storage, or global communications. In addition, information handling systems may include a variety of hardware and software components that may be configured to process, store, and communicate information and may include one or more computer systems, data storage systems, and networking systems.\nInformation handling system performance has steadily increased over time as processing components have increased their ability to process greater amounts of information in reduced amounts of time. Another factor to improved information handling system performance is the greater speeds at which buses within information handling systems are able to communication information between processing components. For example, the recently introduced PCI Express bus communicates information at rates that are considerably more rapid than its predecessor PCI bus. Generally, the more rapid processing and communication of information occurs due to the increased number of transistor devices integrated within processing components to perform processing operations and the greater number of operating cycles that the processing components run per operating time period. However, the ability of processing components within an information handling system to more rapidly process information typically means that the amount of power used by the information handling system also increases. For portable information handling systems that operate on internal battery power having limited duration, the increased power consumption often translates into reduced battery life for a give battery charge.\nOne method for reducing power consumption of an information handling system is to place processing components in an idle state when the processing components are inactive. As an example, the PCI Express standard defines an Active State Power Management (ASPM) protocol in which transmission of characters over a PCI Express link are stopped during idle periods in order to reduce power during times of inactivity. The ASPM protocol states that an inactivity timer set at a value determinable by the processing component manufacturer determines when the communication link should transition to a low power state. Before the transition to the low power state, the PCI Express link transmits characters during times of inactivity in order to maintain the serial link synchronization, such as the synchronization of a phased lock loop (PLL). In the idle state, link synchronization is sacrificed for power savings and re-established once the link becomes active again. However, re-establishing synchronization of a PCI Express link takes some time and thus introduces some delay in the communication of information across the link. An ideal inactivity timer value for assertion of low power states at a PCI Express link depends upon link traffic patterns which vary widely depending upon the type of components involved, system activity, the application being run or other factors. Component manufacturers are motivated to select an inactivity timer value that minimizes the impact on component performance or is optimized for a give set of conditions that may differ substantially from actual operating conditions."}
{"text": "The invention relates to a video circuit for generating video signals for displaying pictures on a picture screen, the video circuit comprising a picture signal source, at least one luminance circuit coupled thereto for generating a luminance signal, and a protection circuit for protecting the picture screen from burning-in due to a prolonged stationary picture.\nVideo signals of this type are used for generating at least a Y-luminance signal and/or R-, G- and B- luminance signals for the pictorial display of information from an information processing device, including Teletext and Viewdata systems as well as video-games.\nThe luminance signals can be applied to corresponding inputs of a display device or can be modulated on a carrier and applied to an aerial input of a television receiving device.\nIt frequently happens that a certain picture remains stationary on the picture screen for a long period of time, so that bright portions of the picture may cause an accelerated local burning-in of the picture screen. Particularly, it often happens that users forget to switch off a television game so that, for example, a line pattern of a playing-field remains on the screen for a whole night. The video circuit is provided with a protection circuit which tries to obviate this harmful effect.\nA video circuit of the type defined in the preamble is known from Canadian patent application No. 267,742. In this video circuit the protection circuit comprises a temperature-dependent element which has the effect that the picture to be displayed is slowly moved over a short distance across the picture screen due to variations in the ambient temperature. However, it appeared that in rooms having a thermostatically controlled heating system, a temperature is adjusted at the beginning of a night which remains very constant thereafter, which is partly caused by the fact that, for example, doors and windows remain closed for a long period of time so that the picture can still remain stationary hours on end in the same position on the picture screen for the remainder of the night."}
{"text": "Field\nThe present disclosure relates, generally, to an apparatus and, more specifically, to an attachable headrest with an eye cover.\nBackground\nA headrest may provide a user with neck and head support during times of rest. In some circumstances, the headrest may be used while the user is sitting upright. For example, the user may be sitting on a seat while traveling (e.g., on a seat while traveling in an airplane or bus). If the user falls asleep while sitting upright, the user's head may deviate from the headrest. For example, the user's head may move to the left, right, or forward such that their head is no longer supported by the headrest. Further, as the user's head moves, left, right, or forward, the headrest may also move because the headrest is no longer secured in place by the weight or force from the user's head. Accordingly, certain problems may exist in existing designs of headrests.\nAn eye cover may block substantial amounts of light while the user is resting. However, existing eye covers may not help to prevent movement of the user's head during times of rest. Accordingly, the problems described supra with respect to movement of the user's head during times of rest may not be overcome by existing eye covers. Therefore, a need exists for an attachable headrest with eye cover that reduces movement of the user's head during times of rest."}
{"text": "This invention relates to heart-lung machines, and more particularly to an automatic occlusion clamp that prevents retrograde blood flow through a centrifugal blood pump.\nIn a typical heart-lung machine of the type used in open-heart surgery, an extracorporeal blood circuit is established from the patient\"\"s venous system through a venous reservoir, a blood pump and an oxygenator to the patient\"\"s arterial system. This circuit takes over the function of the patient\"\"s heart and lungs while the patient\"\"s heart is stopped for the surgery.\nThe blood pump in the extracorporeal circuit is typically either a roller pump or a centrifugal pump. An advantage of the roller pump is that the blood line is always occluded at some point in the pump, so that there can be no retrograde blood flow in the circuit when the pump is stopped. A disadvantage is that a roller pump can convey air boluses through the line. Such air boluses need to be dealt with separately in order to avoid serious injury or death of the patient.\nCentrifugal pumps, on the other hand, do not convey air boluses through the circuit. If an air bolus is applied to the intake of a centrifugal pump, the pump deprimes and needs to be reprimed and restarted. Because of this safety factor, centrifugal pumps are the preferred choice of many perfusionists.\nA problem with centrifugal pumps is that they do not occlude the line, and that consequently, retrograde blood flow through the pump is possible when the pump is stopped, due to the pressure head in the patient\"\"s arterial system. Such a retrograde flow drains blood from the patient and is medically unacceptable.\nRetrograde flow through a centrifugal pump has been averted in the prior art by inserting into the extracorporeal circuit a one-way duckbill valve. The use of such a valve, though effective, does have two disadvantages: first, the shape of a duckbill valve causes turbulence at the exit end of the valve and thereby causes some hemolysis; and second, the duckbill valve, because it is part of the extracorporeal circuit, is a single-use item and therefore adds to the cost of patient disposables.\nThe present invention overcomes the above-described deficiencies of the prior art by providing an electrically or pneumatically operated occlusion clamp external to the blood line that is arranged to occlude the line, preferably between the centrifugal pump and the oxygenator, whenever the forward blood flow in the line essentially ceases. The clamp of this invention is designed to operate in a fail-safe mode, i.e. it will clamp the line shut in the event of a power failure. In combination with conventional devices that shut down the pump in the event of a low blood level condition in the venous reservoir and/or of the detection of an air bolus, the present invention provides a triple-action safety system."}
{"text": "Multiprocessing computer systems include two or more microprocessors that may be employed to perform computing tasks. A particular computing task may be performed on one microprocessor while other microprocessors perform unrelated computing tasks. Alternatively, components of a particular computing task may be distributed among multiple microprocessors to decrease the time required to perform the computing task as a whole.\nA popular architecture in commercial multiprocessing computer systems is the symmetric multiprocessor (SMP) architecture. Typically, an SMP computer system comprises multiple microprocessors connected through a cache hierarchy to a shared bus. Additionally connected to the bus is a memory, which is shared among the microprocessors in the system. Access to any particular memory location within the memory occurs in a similar amount of time as access to any other particular memory location. Since each location in the memory may be accessed in a uniform manner, this structure is often referred to as a uniform memory architecture (UMA).\nProcessors are often configured with internal caches, and one or more caches are typically included in the cache hierarchy between the microprocessors and the shared bus in an SMP computer system. Multiple copies of data residing at a particular main memory address may be stored in these caches. In order to maintain the shared memory model, in which a particular address stores exactly one data value at any given time, shared bus computer systems employ cache coherency. Generally speaking, an operation is coherent if the effects of the operation upon data stored at a particular memory address are reflected in each copy of the data within the cache hierarchy. For example, when data stored at a particular memory address is updated, the update may be supplied to the caches that are storing copies of the previous data. Alternatively, the copies of the previous data may be invalidated in the caches such that a subsequent access to the particular memory address causes the updated copy to be transferred from main memory. For shared bus systems, a snoop bus protocol is typically employed. Each coherent transaction performed upon the shared bus is examined (or “snooped”) against data in the caches. If a copy of the affected data is found, the state of the cache line containing the data may be updated in response to the coherent transaction.\nUnfortunately, shared bus architectures suffer from several drawbacks which limit their usefulness in multiprocessing computer systems. As additional microprocessors are attached to the bus, the bandwidth required to supply the microprocessors with data and instructions may exceed the peak bandwidth of the bus. Thus, some microprocessors may be forced to wait for available bus bandwidth and the performance of the computer system will suffer when the bandwidth requirements of the microprocessors exceed available bus bandwidth.\nAdditionally, adding more microprocessors to a shared bus increases the capacitive loading on the bus and may even cause the physical length of the bus to be increased. The increased capacitive loading and extended bus length increases the delay in propagating a signal across the bus. Due to the increased propagation delay, transactions may take longer to perform. Therefore, the peak bandwidth of the bus may decrease as more microprocessors are added.\nA common way to address the problems incurred as more microprocessors and devices are added to a shared bus system, is to have a hierarchy of buses. In a hierarchical shared bus system, the microprocessors and other bus devices are divided among several low-level buses. These low-level buses are connected by high-level buses. Transactions are originated on a low-level bus, transmitted to the high-level bus, and then driven back down to all the low level-buses by repeaters. Thus, all the bus devices see the transaction at the same time and transactions remain ordered. The hierarchical shared bus logically appears as one large shared bus to all the devices. Additionally, the hierarchical structure overcomes the electrical constraints of a single large shared bus.\nCo-Pending U.S. patent application Ser. No. 09/815,442 entitled “Method and Apparatus for Efficiently Broadcasting Transactions between a First Address Repeater and a Second Address Repeater” discloses a novel architecture that includes a high-level bus, a plurality of low-level buses, and a novel distributed arbiter. As the efficiency of a computer system that includes the above architecture is dependent upon the proper operation of the distributed arbiter, a need exists for methods of verifying the consistency of the distributed arbiter."}
{"text": "LED packaging technique is one of the focuses in the development of the semiconductor industry. Packaging technique for lateral LEDs and vertical LEDs by wire bonding, and packaging technique for flip-chip LEDs by flip-chip bonding are all being studied for improving brightness and reliability as well as reducing overall volume of the packaged LEDs.\nIn terms of packaging technique by flip-chip bonding, chip-scale packaged LED device has been developed for the purpose of miniaturization. However, there is still room for improvement regarding brightness, volume, reliability, yield, etc."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrical connector assembly, and more particularly to an electrical connector assembly for electrically connecting to a printed circuit board.\n2. Description of Related Art\nAn electrical connector assembly 100′ shown in FIG. 1 comprises a housing 5′, a plurality of terminals received in the housing 5′, a chip module 9′ mounted in the housing 5′, a stiffener 6′ surrounding the housing 5′, a metal clip 7′ attached on an end of the stiffener 6′ for pressing the chip module 9′ and a lever 8′ engaging with the other end of the stiffener 6′ for locking the metal clip 7′. The metal clip 7′ defines a pressing portion 71′. When the metal clip 7′ is pressed downwardly, the pressing portion 71′ abuts against the chip module 9′ into the housing 5′.\nHowever, the chip module 9′ is directly mounted between the housing 5′ and the metal clip 7′. When the lever 8′ is rotated to press down the chip module 9′, the chip module 9′ bears two kinds of forces applied by the terminals and the metal clip 7′. In this case, the terminals received in the housing 5′ generate a nonuniform deformation due to an uneven force applied on, which can lead to an open electrical connection between the chip module 9′ and the terminals in a slight deformation section and destroy the terminals in a greater deformation section. So the chip module 9′ can be crashed by the pressing portion 71′.\nAn improved electrical connector assembly 200′ shown in FIG. 2 comprises a housing 10′, a plurality of terminals received in the housing 10′, a chip module 14′ mounted in the housing 10′, a stiffener 11′ surrounding the housing 10′, a metal clip 12′ and a lever 13′ engaging with two opposite ends of the metal clip 12′ and a heat sink 15′ attached on the chip module 14′. The heat sink 15′ defines a flange 151′. The metal clip 12′ presses the chip module 14′ via the flange 151′, which improves the pressing character thereof and electrical connection efficiency between the chip module 14′ and the terminals received in the housing 10′. However, the heat sink 15′ is mainly used for dissipating heat radiated by the chip module 14′ and the material of the heat sink 15′ must have a good heat dissipation, hence leading to a higher cost.\nThus, there is a need to provide a new electrical connector assembly that overcomes the above-mentioned problem."}
{"text": "1. Field of the Invention\nThe present invention is generally related to a safety device for a marine vessel and, more particularly, to a detector that senses the presence of a human being in the region of the marine propulsion device and deactivates the device to prevent injury to the human being.\n2. Description of the Prior Art\nMany different types of marine propulsion devices are well known to those skilled in the art. For example, outboard motors that are attached to the transom of a marine vessel, stern drive systems that extend in a rearward direction from the transom of a marine vessel, docking thrusters, and bow thrusters are well known to those skilled in the art. In addition to bow thrusters, certain types of docking thruster systems used in conjunction with marine vessels incorporate a plurality of propulsors that are responsive to the joystick manipulations of a marine vessel operator.\nMany types of infrared sensors are well known to those skilled in the art. These types of sensors have been used in many applications where the presence of a human being is detected and the detection is used to activate or deactivate a device, such as a lighting system.\nU.S. Pat. No. 6,142,841, which issued to Alexander et al on Nov. 7, 2000, discloses a water jet docking control system for a marine vessel. Several versions of the docking system are disclosed in this patent. Once system utilizes pressurized liquid at three or more positions of a marine vessel in order to selectively create thrust that moves the marine vessel into desired locations and according to chosen movements. Electrical embodiments of the system utilize one or more pairs of impellers to cause fluid to flow through the outlet conduits in order to provide thrust on the marine vessel. In one embodiment of the device, a cross thrust conduit is associated with the marine vessel to direct fluid flow in a direction perpendicular to a centerline of the marine vessel and a pair of outlet conduits are associated with the marine vessel to direct flows of fluid in directions which are neither parallel nor perpendicular to a centerline of the marine vessel.\nU.S. Pat. No. 5,283,427, which issued to Phillips et al on Feb. 1, 1994, describes a night sight for a missile launcher comprising an image intensifier tube, a reticle, and an objective lens. The missile launcher night sight has an objective lens with a field of view of at least 22 degrees. The output image of the objective lens is intensified by a variable gain light intensifier tube and the output of the intensifier is viewed through an eyepiece. A reticle pattern etched on a glass substrate and filled with titanium dioxide is illuminated by adjustable brightness LED\"\"s positioned at points on the periphery of the substrate. The reticle is disposed between the light intensifier and the eyepiece and substantially duplicates the function of reticles used in daysights.\nU.S. Pat. No. 6,160,948, which issued to McGaffigan on Dec. 12, 2000, describes optical light pipes with laser light appearance. A simulated laser light system includes a light source which emits substantially parallel light rays and a prismatic element having a plurality of prismatic surfaces. The prismatic surfaces define a curved surface either on the inside or the outside of the prismatic element. The prismatic surfaces redirect the light rays emitted from the prismatic element in a plurality of plants which are perpendicular to a tangent to the curved surface. The simulated laser light system appears to emit light from a centerline of the prismatic element when the light is actually emitted from an exterior surface of the prismatic element.\nU.S. Pat. No. 3,950,645, which issued to Rotstein et al on Apr. 13, 1976, describes an infrared detection tube. An imaging device comprises a cooling means for lowering the temperature of the device, a wafer of preselected infrared sensitive material, a layer of thin infrared transparent electrically conductive material deposited on one side of the wafer, a wire mosaic, the wire mosaic being in intimate contact with the remaining side of the wafer, opaque material filling a preselected area between the wires of the wire mosaic, and a silver paste interposed between the wafer and the tips of the wires of the wire mosaic. The resistance of the infrared sensitive wafer cooperates with a grid current flow and reduces the bias permitting flow of electrons according to infrared radiation impinging on the wafer.\nU.S. Pat. No. 3,936,822, which issued to Hirschberg on Feb. 3, 1976, describes a method and apparatus for detecting weapon fire. The method and apparatus disclosed in this patent relates to the automatic detection of the firing of weapons, such as small arms, or the like. Radiant and acoustic energy produced upon occurrence of the firing of the weapon and emanating from the muzzle thereof are detected at known, substantially fixed, distances therefore. Directionally sensitive radiant and acoustic energy transducer means are directed toward the muzzle to receive the radiation and acoustic pressure waves therefrom may be located adjacent each other for convenience.\nU.S. Pat. No. 5,018,872, which issued to Suszynski et al on May 28, 1991, describes a probe assembly for an infrared thermometer. The probe assembly is adapted for insertion into a patient\"\"s ear canal and facilitates a rapid measurement of the patient\"\"s body temperature with very high accuracy. The probe assembly includes a plastic outer tube sized to fit snugly in the patient\"\"s ear canal, with an elongated heat sink and an infrared sensor located within the tube.\nU.S. Patent No. 6,100,803, which issued to Chang on Aug. 8, 2000, describes an infrared illuminative warning detector. The detector includes a base seat formed with at least four perforations for two light shades and two detector heads to insert in. A bulb in installed in each light shade. An infrared detector is disposed in each detector head for detecting alien article within a detection range and lighting up the bulb. Each light shade and detector head is disposed with at least one shifting mechanism for freely changing operation position. The light shades and the detector heads on the base seat are able to detect in different directions at the same time to provide a warning and illumination function in the case of intrusion of alien article.\nU.S. Pat. No. 5,987,205, which issued to Moseley et al on Nov. 16, 1999, describes an infrared energy transmissive member and radiation receiver. The infrared energy transmissive member is used for conducing infrared energy from a first end of the infrared energy transmissive member to a second end disposed adjacent an infrared responsive circuit component of an infrared receiver. The member comprises a flexible hollow plastic tube. Preferably, the plastic tube comprises an ester based polyurethane tube and has a malleable rod disposed in the hollow plastic tube to allow the hollow plastic tube to be bent into a claimed configuration. The infrared responsive circuit component may control a fluorescent lamp dimming circuit or a window treatment control circuit, thereby allowing remote control of the lamp or window treatment.\nThe patents described above are hereby expressly incorporated by reference in the description of the present invention.\nA safety device for a marine vessel, made in accordance with the present invention, comprises an electrical device, such as an electrical marine propulsion device, and an infrared sensor having an output signal representing a change in heat intensity within a viewing angle of the infrared sensor. It should be understood that the electrical device can be an electrically driven propulsor, such as a bow thruster for a marine vessel, or any other marine propulsion device that has an electrical element that can be deactivated for the purpose of deactivating the marine propulsion device. As an example, a conventional outboard motor incorporating an ignition system which is electrical is an electrical device within the meaning of the description of the present invention. The outboard motor can be deactivated, for example, by deactivating the ignition system.\nThe present invention further comprises a tube having a central cavity, the central cavity having a first end and a second end. The infrared sensor is disposed at the first end of the central cavity to receive infrared signals passing along a line of sight defined by the tube in a direction through the central cavity from the second end of the central cavity toward the first end. The viewing angle of the infrared sensor is defined by a length of the tube and a diameter of the tube. The tube is attached to the marine vessel to direct the line of sight in a generally downward direction toward a surface of water proximate the marine vessel to detect a change in heat intensity in a region proximate the surface of water and proximate the marine vessel. A preferred embodiment of the present invention further comprises an alarm circuit that is connected in signal communication with the infrared sensor and an electrical communication with the electrical device. The alarm circuit is configured to deactivate the electrical device in response to receipt of the output signal from the infrared sensor.\nThe magnitude of the length of the tube in a preferred embodiment of the present invention, is at least three times the magnitude of the diameter of the tube. More specifically, one embodiment of the present invention that has been empirically tested comprises a tube whose length is four times the magnitude of its diameter.\nThe alarm circuit can have an audible alarm output that is sounded in response to receipt of the output signal from the infrared sensor. The electrical device can be an electric motor with a propulsor attached to an output shaft of the electric motor or, as described above, it can be any marine propulsion device or system that incorporates an electric component that can be affected for the purpose of deactivating the electrical device.\nIn a particularly preferred embodiment of the present invention, the tube is opaque and the central cavity is circular in cross section. The infrared sensor in a preferred embodiment of the present invention comprises two infrared sensing elements that are associated with each other in a differential arrangement."}
{"text": "In the presence of catalysts, curable organosilicone compositions undergo hydrosilation, a reaction involving the addition of a siliocon-hydrogen bond across a pair of aliphatic carbon atoms linked by multiple bonds. Catalysts used for reactions of this type are metals, most notably platinum, rhodium, iridium, and palladium, and compounds thereof. Hydrosilation has found widespread use in the production of silicone materials and organosilanes. Compared with platinum, relatively little work has been done on rhodium- and iridium-containing hydrosilation catalysts. Various chloro-rhodium compounds, wherein the valence of rhodium (Rh) is 1, particularly [RhCl(CO).sub.2 ].sub.2 and RhCl(CO)[P(C.sub.6 H.sub.5).sub.3 ].sub.3 (A. J. Chalk, J. Organometal. Chem; 21, 207-213 (1970)), [RhCl(C.sub.2 H.sub.4).sub.2 ].sub.2 (U.S. Pat. No. 3,296,291), and RhCl[P(C.sub.6 H.sub.5).sub.3 ].sub.3 (U.S. Pat. No. 3,546,266), are disclosed in the literature to be hydrosilation catalysts. The mono- and bi-metallic pyrazine compounds, (pyrazine)Rh(CO).sub.2 Cl and (pyrazine)Rh.sub.2 (CO).sub.4 Cl.sub.2, have been disclosed (A. L. Balch et al., J. Organometal. Chem., 169, 97 (1979), but no use is specified. R. N. Haszeldine, R. V. Parish and D. J. Perry, J. Chem. Soc. (A), 683 (1969), found that [(C.sub.6 H.sub.5).sub.3 P].sub.2 Rh(CO)Cl was active as a hydrosilation catalyst but that the iridium analogue was inactive."}
{"text": "The hydroelectric power plant of the present invention employs a generating system similar to that employed in my U.S. Pat. No. 4,075,838. As in the patent, there are casing means defining a first chamber filled with a first column of water and a standpipe relatively inclined to the horizontal thereabove to pressurize the first column of water. There is also a housing enclosed in the first chamber which defines a second chamber therein for containing water in a second column horizontally offset from the first column of water. The housing has a series of penstocks about the perimeter thereof which open into the first column of water at a first level above the second chamber, and discharge into the second chamber at a second level below the first level. In addition, there are hydroelectric power generation means including turbines in the penstocks and means for discharging the tailwater from the respective turbines. The tailwater discharge means include elongated hollow-bore, open ended nipples extending into the first chamber from the housing at a third level below the first and second levels and having open ended thimble-like vessels telescopically engaged on the end portions thereof which are relatively remote from the housing, to be extended and retracted in relation to the respective nipples, axially thereof. There are also drive means operable to extend and retract the respective vessels, and float operated check valves in the vessels adjacent the end openings of the vessels relatively remote from the housing, each of which is adapted to permit flow through the bore of the respective nipple in the direction relatively away from the housing toward the first column of water in the first chamber when open, and prevent flow in the opposite direction when closed. As in the patent, control means are connected to the drive means to extend and retract the respective vessels in staggered sequence to one another about the circumference of the housing, with each vessel undergoing an at rest period between the respective retraction and extension stages thereof, when the vessel is fully retracted."}
{"text": "miRNA is a kind of 20-25-nucleotide endogenous small RNA present in cells. It is derived from a DNA that produces no protein, and is generated from hairpin-shaped transcripts. miRNA functions as a repressor that binds to the complementary sequence of 3′-UTR of its target mRNA to induce the inhibition of translation or unstabilization of the target mRNA, and ultimately repress the protein synthesis of the target mRNA. It is known that one miRNA targets several mRNAs and mRNA may also be controlled by several miRNAs.\nmiRNA is currently known to play a very important role in a number of biological processes including regulation of developmental timing, apoptosis, fat metabolism, and hematopoietic cell differentiation among others, and thus has received a great deal of attention in the field of life science. However, while studies on the function of miRNA in the processes of development of diseases such as cancer have considerably progressed, studies on the role of miRNA in the skin science are relatively insufficient. Recently, it was reported that miR-203 targets p63 mRNA in keratinocytes to inhibit the differentiation of the keratinocytes and that miR-125b inhibits the differentiation of skin stem cells and inhibits the division of keratinocytes, but there are little or no reports of skin-related specific miRNA and the function thereof. In addition to miR-203 and miR-125b, several specific miRNAs are sufficiently expected to play a biologically important role in the skin. In addition, miRNA can be used in various applications, including biological markers and materials, and thus it is concluded that studies on the function of miRNA in the skin can be applied in various fields, including the cosmetic and medical fields.\nThe skin is roughly divided into epidermis and dermis, and melanocytes present in epidermis synthesize a pigment (melanin) that shows the inherent color of the skin. Biological factors that control skin pigments can be roughly divided into factors that make melanosome in melanocytes, and factors that transport melanosome. Of the factors that make melanosome, the most well-known factors include tyrosinase (TYR) and the like, and factors known to be involved in melanosome transport include Rab27a and the like. Recently, it was known that transcription factors such as PAX3, MITF and the like also participate in the pigment regulation process. Thus, these factors that regulate the pigment become important targets for skin whitening, black hair production or gray hair prevention."}
{"text": "The invention relates to power transmission systems and components of the type including a pulley having a circumferentially spaced portion of its drive face spaced unequal distances from the pulley axis of rotation, but more particularly, the invention relates to a tension compensating gear set of the eccentric type with a gear member having external teeth that mesh with internal teeth of a floating, annular gear member.\nEccentric type gear sets are used for a variety of purposes in power transmission systems with a flexible power transmission member. For example, such gear sets may be used to increase the minimum bend radius of a flexible power transmission member such as illustrated in German Patent No. 115,719; as a speed reduction means as shown in U.S. Pat. No. 3,190,149; or as a means for compensating for differences in tension in a flexible power transmission member such as shown in U.S. Pat. No. 4,501,576. The term \"flexible power transmission member\" as used herein includes a toothed or synchronous belt, a V-belt, a V-ribbed belt, a flat belt, and a chain.\nRadial disengagement of meshing teeth of a tension compensating gear set is characteristically inhibited by a flexible power transmission member, and lateral disengagement of meshing teeth of the gear set may be prevented by lateral side plates of the type as disclosed in aforementioned U.S. Pat. No. 4,501,576 or German OS 31 00 922.\nProblems associated with such eccentric type gear sets as used with a flexible power transmission member, are that of noise and vibration as caused by the annular gear member not having a fixed or precise center of rotation. Variations in tension of the flexible power transmission member cause a radially inward force to be applied to the annular gear member causing its gear teeth to wedge with the teeth of the gear member and carry radial loads; this radial force is in addition to the tangential force carried by the gear teeth at their meshing flank surfaces."}
{"text": "Opto-electronic devices that make use of organic materials are becoming increasingly desirable for a number of reasons. Many of the materials used to make such devices are relatively inexpensive, so organic opto-electronic devices have the potential for cost advantages over inorganic devices. In addition, the inherent properties of organic materials, such as their flexibility, may make them well suited for particular applications such as fabrication on a flexible substrate. Examples of organic opto-electronic devices include organic light emitting devices (OLEDs), organic phototransistors, organic photovoltaic cells, and organic photodetectors. For OLEDs, the organic materials may have performance advantages over conventional materials. For example, the wavelength at which an organic emissive layer emits light may generally be readily tuned with appropriate dopants.\nOLEDs make use of thin organic films that emit light when voltage is applied across the device. OLEDs are becoming an increasingly interesting technology for use in applications such as flat panel displays, illumination, and backlighting. Several OLED materials and configurations are described in U.S. Pat. Nos. 5,844,363, 6,303,238, and 5,707,745, which are incorporated herein by reference in their entirety.\nOne application for phosphorescent emissive molecules is a full color display. Industry standards for such a display call for pixels adapted to emit particular colors, referred to as “saturated” colors. In particular, these standards call for saturated red, green, and blue pixels. Color may be measured using CIE coordinates, which are well known to the art.\nOne example of a green emissive molecule is tris(2-phenylpyridine) iridium, denoted Ir(ppy)3, which has the following structure:\n\nIn this, and later figures herein, we depict the dative bond from nitrogen to metal (here, Ir) as a straight line.\nAs used herein, the term “organic” includes polymeric materials as well as small molecule organic materials that may be used to fabricate organic opto-electronic devices. “Small molecule” refers to any organic material that is not a polymer, and “small molecules” may actually be quite large. Small molecules may include repeat units in some circumstances. For example, using a long chain alkyl group as a substituent does not remove a molecule from the “small molecule” class. Small molecules may also be incorporated into polymers, for example as a pendent group on a polymer backbone or as a part of the backbone. Small molecules may also serve as the core moiety of a dendrimer, which consists of a series of chemical shells built on the core moiety. The core moiety of a dendrimer may be a fluorescent or phosphorescent small molecule emitter. A dendrimer may be a “small molecule,” and it is believed that all dendrimers currently used in the field of OLEDs are small molecules.\nAs used herein, “top” means furthest away from the substrate, while “bottom” means closest to the substrate. Where a first layer is described as “disposed over” a second layer, the first layer is disposed further away from substrate. There may be other layers between the first and second layer, unless it is specified that the first layer is “in contact with” the second layer. For example, a cathode may be described as “disposed over” an anode, even though there are various organic layers in between.\nAs used herein, “solution processible” means capable of being dissolved, dispersed, or transported in and/or deposited from a liquid medium, either in solution or suspension form.\nA ligand may be referred to as “photoactive” when it is believed that the ligand directly contributes to the photoactive properties of an emissive material. A ligand may be referred to as “ancillary” when it is believed that the ligand does not contribute to the photoactive properties of an emissive material, although an ancillary ligand may alter the properties of a photoactive ligand.\nAs used herein, and as would be generally understood by one skilled in the art, a first “Highest Occupied Molecular Orbital” (HOMO) or “Lowest Unoccupied Molecular Orbital” (LUMO) energy level is “greater than” or “higher than” a second HOMO or LUMO energy level if the first energy level is closer to the vacuum energy level. Since ionization potentials (IP) are measured as a negative energy relative to a vacuum level, a higher HOMO energy level corresponds to an IP having a smaller absolute value (an IP that is less negative). Similarly, a higher LUMO energy level corresponds to an electron affinity (EA) having a smaller absolute value (an EA that is less negative). On a conventional energy level diagram, with the vacuum level at the top, the LUMO energy level of a material is higher than the HOMO energy level of the same material. A “higher” HOMO or LUMO energy level appears closer to the top of such a diagram than a “lower” HOMO or LUMO energy level.\nAs used herein, and as would be generally understood by one skilled in the art, a first work function is “greater than” or “higher than” a second work function if the first work function has a higher absolute value. Because work functions are generally measured as negative numbers relative to vacuum level, this means that a “higher” work function is more negative. On a conventional energy level diagram, with the vacuum level at the top, a “higher” work function is illustrated as further away from the vacuum level in the downward direction. Thus, the definitions of HOMO and LUMO energy levels follow a different convention than work functions.\nMore details on OLEDs, and the definitions described above, can be found in U.S. Pat. No. 7,279,704, which is incorporated herein by reference in its entirety."}
{"text": "Field\nEmbodiments relate generally to sodium free glasses and more particularly to fusion formable sodium free glasses which may be useful in photochromic, electrochromic, Organic Light Emitting Diode (OLED) lighting, or photovoltaic applications, for example, thin film photovoltaics.\nTechnical Background\nThe fusion forming process typically produces flat glass with optimal surface and geometric characteristics useful for many electronics applications, for instance, substrates used in electronics applications, for example, display glass for LCD televisions.\nOver the last 10 years, Corning fusion glass products include 1737F™, 1737G™, Eagle2000F™, EagleXG™, Jade™, and Codes 1317 and 2317 (Gorilla Glass™). Efficient melting is generally believed to occur at a temperature corresponding to a melt viscosity of about 200 poise (p). These glasses share in common 200p temperatures in excess of 1600° C., which can translate to accelerated tank and electrode corrosion, greater challenges for fining due to still more elevated finer temperatures, and/or reduced platinum system life time, particularly around the finer. Many have temperatures at 3000 poise in excess of about 1300° C., and since this is a typical viscosity for an optical stirrer, the high temperatures at this viscosity can translate to excessive stirrer wear and elevated levels of platinum defects in the body of the glass.\nMany of the above described glasses have delivery temperatures in excess of 1200° C., and this can contribute to creep of isopipe refractory materials, particularly for large sheet sizes.\nThese attributes combine so as to limit flow (because of slow melt rates), to accelerate asset deterioration, to force rebuilds on timescales much shorter than product lifetimes, to force unacceptable (arsenic), expensive (capsule) or unwieldy (vacuum fining) solutions to defect elimination, and thus contribute in significant ways to the cost of manufacturing glass.\nIn applications in which rather thick, comparatively low-cost glass with less extreme properties is required, these glasses are not only overkill, but prohibitively expensive to manufacture. This is particularly true when the competitive materials are made by the float process, a very good process for producing low cost glass with rather conventional properties. In applications that are cost sensitive, such as large-area photovoltaic panels and OLED lighting, this cost differential is so large as to make the price point of LCD-type glasses unacceptable.\nTo reduce such costs, it is advantageous to drive down the largest overall contributors (outside of finishing), and many of these track directly with the temperatures used in the melting and forming process. Therefore, there is a need for a glass that melts at a lower temperature than those aforementioned glasses.\nFurther, it would be advantageous to have a glass useful for low temperature applications, for instance, photovoltaic and OLED light applications. Further, it would be advantageous to have a glass whose processing temperatures were low enough that the manufacturing of the glass would not excessively consume the energy that these applications are aiming to save."}
{"text": "The present invention relates to swimming apparatus, in particular, to flippers to be worn on the feet of the swimmer to enhance mobility.\nWhenever man has endeavored to develop in an environment in which he was alien, he had to look for inspiration and teaching to those organs which nature gave to the creatures native to that environment, which organs man himself lacked. For example, the ailerons of an airplane are related to the wings of the birds, and the flippers which divers and swimmers employ to supplement the fin-like structure of hands and feet are similar to the webbing and the fins of fish and batrachians (frogs, toads and the like).\nIt is evident that it is the movement of the water over the surface of a fin that creates the push which makes the swimmer (fish or man) proceed with increased mobility. However, this same surface creates friction, turbulence and cavitation within the water producing a diminished result in the applied force of the swimmer. Because the natural fins of man (hands or feet) are relatively small and ineffective, man has begun to utilize supplemental flippers which he fastens to his feet which enable these lower limbs to better move and flex. In general, most flippers used are formed of a single membrane or sheet having an upper and lower surface only. Recently, however, some flippers have been made utilizing a second membrane, so that a kind of a flat, hollow compartment is formed between the two membranes in the largest part of the flipper. This arrangement permits the flipper to hold back water on one side, while permitting its replacement on the other side. This flow of water resembles the flowing of water which one observes within the locomotor tubes of certain cephalopodes and molluscs, particularly the cuttlefish. As a result, a slight improvement in the effect of the applied force of the fin is observed. However, such construction does not enable the expansion or flow of water outwardly of the fins in both the downward and in the upward stroke of the swimmer, similar to that experienced by the batrachians.\nOne may liken this latter effect to that experienced by a bicycle rider while peddling. The propelling force of the bicycle is limited when the rider only pushes downwardly on the pedals alternately with his feet. However, by utilizing the chain drive which maintains the pedal in constant motion, the efforts of raising the opposing feet once the pedals have passed a dead center is lessened and an increased force efficiency is obtained. This is easily seen merely by fastening each foot to the pedal near the metatarsus. This method of pedalling the bicycle is similar to the push-pull work used to operate a crank.\nWith the foregoing in mind, the object of the present invention is to improve the construction of the flippers, used by swimmers, so that they are more effective to the swimmer in both the upward and downward strokes.\nIt is a further object of the present invention to provide an improved flipper which enables the flow of water through it, without creating turbulence, or cavitation, in the water.\nIt is another object of the present invention to provide a flipper in which greater force efficiency and a synergistic effect is obtained in both the downward and upward strokes.\nThe foregoing objects, together with other advantages and objects, will be apparent from the following disclosure of the present invention."}
{"text": "(a) Field of the Invention\nThe present invention relates to a circuit system, and more particularly to a variable resistance linear limit voltage circuit that may, by way of example and not limitation, be used to match the rated saturation voltage VS of different rechargeable cells during charging, and that employs light emitting diodes as both illumination and impedance devices so as to reduce heat generated by the impedance devices while still achieving the desired limit voltage control.\n(b) Description of the Prior Art\nU.S. Pat. No. 5,118,993 and Europe Patent No. 0487204, both granted to the applicant of the present invention, disclose a multi-voltage output circuit that uses a positive voltage drop from a diode or a zener voltage from a zener diode to control an output voltage supplied to rechargeable cells directly connected in parallel at the output of the multi-voltage output circuit during a charging process, or the positive voltage drop from multiple diodes directly connected in parallel with the cells to provide a limit-voltage-divided charging current. In either application, the circuit provides a regulated voltage V0 when the terminal voltage of the cells accumulates along with the charging current and rises up to such extent close to, and eventually becomes identical with the positive voltage drop value. However, the positive voltage drop of the diodes in such an arrangement has a gradient of approximately 0.7V difference depending on the number of diodes connected in series, and therefore it is very difficult to match the rated saturation voltage VS of the cells by changing the number of diodes connected in series when the positive voltage drop of the diodes connected in parallel is not of the same value as that of the VS. Therefore, diodes must be directly connected in parallel with the cells being charged, and such connection creates the following defects:\n1. In the absence of additional connection in series of a proper limiting current, the charging current IB decreases when a composite regulated voltage V0 is generated by the positive voltage drop of the diode and the terminal voltage of the cells. As a result, the diodes are vulnerable to being burnt out due to the significantly increased current passing through the diodes, as illustrated in FIG. 1 of the accompanying drawings of the present invention; and\n2. The positive voltage drop value of the diodes is not consistent with the rated saturation voltage VS required by the cells. If the value is lower than VS, the charging current IB passing through the cells gets too small and consequently, the charging process becomes too slow or the charging current is insufficient, as illustrated in FIG. 2. On the other hand, if the value gets higher than VS, the cells will be overcharged.\nThe primary purpose of the present invention is to provide a linearly variable resistance limit voltage circuit system including a light emitting display.\nThe limit voltage circuit of the invention may be used in a variety of applications, one of which is the application of matching the saturation voltage of a rechargeable cell to ensure fully saturated charging, while at the same time preventing damage to the cell due to overcharging, reducing thermal loss from the limit voltage circuit, and providing a light emitting display when required."}
{"text": "1. Field of the Invention\nThe present application relates to a lubrication system for a two-cycle engine and, more particularly, to a lubrication system that incorporates a lubrication pump that pressurizes and delivers lubricant to a portion of a two-cycle engine.\n2. Description of Related Art\nIn all fields of engine design, there is an increasing emphasis on obtaining more effective emission control. Recent two-cycle engines, therefore, incorporate a lubricant pump to deliver a desired amount of lubricant to lubricate internal portions of the engines. Mechanically operated pumps can be used as the lubricant pump. Such mechanical pumps, however, are not easily controlled to provide highly precise amounts of lubricant in response to engine operations. Electrically operable pumps tend to replace the mechanical pumps because higher precision controls are more widely available with such electrical pumps.\nThe electrical pumps can periodically pressurize lubricant under the control of a control device, such as, for example, an electronic control unit (ECU). The ECU can control a frequency of the periodic pressurization with, for example, an electronic control signal configured to operate the pump in accordance with a desired duty cycle. The higher the frequency, the greater the amount of the lubricant.\nAn electromagnetic solenoid pump is one type of such electrical pump. Japanese Laid Open Patent Publication 10-37730 discloses a lubrication system incorporating such an electromagnetic solenoid pump. The solenoid pump has a pumping piston reciprocally disposed in a pump housing. A plunger is coupled with the pumping piston. An electromagnetic solenoid can actuate the plunger. A control device controls the solenoid to selectively actuate or release the plunger such that the pumping piston periodically pressurizes the lubricant.\nThe control device disclosed in Japanese Laid Open Patent Publication 10-37730 has a control map that provides an amount of lubricant required by the engine versus an engine speed and determines a frequency of energization of the solenoid using the control map. The solenoid pump thus can pressurize a proper amount of lubricant in response to the engine speed of the engine."}
{"text": "The present invention relates to a methods and systems for limiting risks associated with a rating supplied by an investment bank relating to a company's earning estimate and/or an investment recommendation (“Risks”).\nThe capital markets in a country such as the United States may be the most sophisticated and regulated in the world. Research departments of financial institutions analyze and arrive at various estimations of corporate performance, typically for companies whose securities are publicly traded. This type of research, analysis and estimation are an important component to the efficient functioning of such markets. Research analysts develop and utilize specific models to arrive at quarterly earnings estimates for a given company or companies which they cover. Typically such earnings estimates are expressed in terms of earnings per share.\nConsensus estimates have been developed which reflect an aggregated view of multiple analysts to arrive at a “consensus” number of range of numbers of expected earnings per share for a particular interval, such as a fiscal quarter. One well known provider of such consensus estimates is, for example, is Thomson “First Call®”.\nRecently, however, the integrity of research provide by financial institutions has been called into question where the financial institutions with research departments either have, or hope to, achieve banking relationships with companies covered by their research departments. Presently, regulatory, governmental, legislative and media investigators, class action attorneys and the public are questioning whether analyst reporting may have been unduly influenced by business or banking relationships.\nThe implications of these questions are beginning to play out in class action litigations, the opinions of commentators and proposed regulations. Arguments have been advanced that analysts provide earnings estimates and other subjective data for reasons that may be influenced by business relationships, such as investment banking activity, with the company. Presently there is no way to determine, based upon a consensus earnings estimate, which estimates have been provided from financial institutions with banking relationships from covered companies and which estimates have not.\nAccordingly, what is needed is a method and system for in which a consensus earnings estimate can be provided in such a way to identify and aggregate earnings per share estimates provided by financial institutions with existing or proposed banking or financial relationships with covered companies and those earning per share estimates provided by financial institutions or researchers who do not. Consensus estimates can then be provided in a meaningful manner devoid of any real or imagined influence arising from such relationships."}
{"text": "The present invention relates to a dry etching apparatus for minutely processing patterns of semiconductor circuits or the like.\nPlasma etching as a way of dry etching has been often carried out in process of manufacturing semiconductor devices in recent years.\nA conventional dry etching apparatus will be depicted hereinbelow. FIG. 4 diagrammatically indicates an example of a reactive ion etching apparatus. The dry etching apparatus in FIG. 4 is provided with an upper electrode 2 and a lower electrode 3 in a vacuum chamber 1 via an insulating member 4. The lower electrode 3 is maintained at a constant temperature by a water-circulating type of temperature control device 5 through the circulation of water in an internal passage 6. A high frequency power source 8 is connected to the lower electrode 3 via an impedance matching circuit 7, bringing about plasma between the electrodes 2, 3.\nEtching gas is introduced into the vacuum chamber 1 from a steel gas cylinder 9 through a valve 10 and a mass flow controller 11, which is discharged out by a vacuum pump 12. The pressure inside the vacuum chamber 1 is kept constant by a pressure controller 13.\nMore specifically, a predetermined etching gas is fed from the steel gas cylinder 9 through the valve 10 and the mass flow controller 11 into the vacuum chamber 1 in a direction of an arrow. When a high frequency power from the high frequency power source 8 is applied to the lower electrode 3, the etching gas is turned into the plasmic state. The etching gas in the plasmic state etches a to-be-etched member 14 disposed on the lower electrode 3.\nThe lower electrode 3 in the above conventional dry etching apparatus is schematically shown in a diagram in FIG. 5. The lower electrode is one disclosed in Japanese Laid-open Patent publication No. 6-204180. A surface of an aluminum base 3a of the lower electrode 3 to which the member 14 is to be attracted by the electrostatic force is coated with a film 3b of aluminum oxide. Meanwhile, the other surfaces of the aluminum base 3a which would come into contact with the etching gas are coated with an alumina ceramic 3c. A passage 6 for cooling medium is formed inside the lower electrode 3 so as to control the temperature of the to-be-etched member 14.\nIn the constitution as above, a semiconductor wafer 14 as the to-be-etched member is placed on the lower electrode 3 as indicated in FIG. 4. A high frequency power is applied between the electrodes 2, 3 while the etching gas is fed. At this time, the surface of the semiconductor wafer 14 is charged with negative electricity as a result of the cathode potential drop. In consequence, the lower electrode 3, specifically, aluminum oxide coating 3b is polarized, whereby an electrostatic attraction force is generated between the lower electrode 3 and the semiconductor wafer 14, thereby attracting the semiconductor wafer 14 to the lower electrode 3. In this 10 manner, while the semiconductor wafer 14 is maintained at a constant temperature close to the temperature of the lower electrode 3, plasma etching is carried out.\nIn the conventional method described hereinabove, the volume resistivity of the aluminum oxide coating 3b of the lower electrode 3 is 1.times.10.sup.12 through 1.times.10.sup.13 .OMEGA.cm in the etching process. Thus, it is difficult to remove the charged electricity from the coating 3b. Therefore, the amount of charge on the surface of the semiconductor wafer 14 is increased as the etching proceeds. As a result, the obtained semiconductor device may be broken subsequent to the electrical discharge as it the coating 3b is increasingly charged, or the electrostatic attraction force remaining even after etching may cause difficulties in transferring the wafer."}
{"text": "Conventionally, a washer fluid spraying apparatus for a vehicle etc. that sprays the washer fluid stored in a tank to a windshield is known. In addition, such a washer fluid spraying apparatus alerts a driver with a warning light, etc., when the amount of remaining washer fluid in the tank decreases. Further, one known washer fluid spraying apparatus has a control apparatus that changes the spraying amount of the washer fluid so as to be different from that in a normal operation by intermittently driving a motor that drives a washer pump when the amount of the remaining washer fluid in the tank is decreased (for example, see Japanese Laid Open Utility Model Patent No. 62-192965). Another known washer fluid spraying apparatus includes a control apparatus that decreases the amount of the sprayed washer fluid according to the level of the remaining washer fluid when the amount of the remaining washer fluid in the tank is decreased (for example, see Japanese Laid Open Patent No. 2005-75218).\nHowever, since in the above-mentioned prior art, the amount of the remaining washer fluid is electrically detected, and the washer pump is controlled according to the remaining amount, it is necessary to have electric parts such as a sensor that detects the amount of the remaining washer fluid, a control circuit that contains an amplifier harness that changes the amount of the sprayed washer pump, and a warning light, thereby increasing manufacturing costs."}
{"text": "This invention relates to hand-portable wand devices, and more specifically to novelty wands which are intended to realistically simulate those sparkling, shimmering, magic-wand devices made popular in various childrens fairytales, such as in Walt Disney's cinema fantasy Fantasia (the Sourcer's Apprentice), and `PeterPan`, etc..\nAdditionally, this invention relates to specialt uses of fiberoptics for novelty items such as is conveyed in U.S. Pat. No. 3,727,044 for a decorative Xmas-tree ornament, and in U.S. Pat. No. 4,186,425 for a type of illuminated jewelry apparel. Plus, this invention employs the optional use of either a battery powered lighting means somewhat akin to that employed by the `StarWars` simulated ray-sword of U.S. Pat. No. 4,208,701, or the luminous-kite toy of U.S. Pat. No. 4,715,564, which employs the use of chemiluminescense `light-sticks` such as are described in the reactant lighting compositions made popular by the Amer. Cyanamide Corp. such as in U.S. Pat. No. 4,698,183 (specifically a hi-output/short-duration formula) which also lists many of their prior (some now of public domain) chemiluminescent based patents.\nHowever, none of the above, nor other patents reviewed by this inventor, show anticipation of a device which could rather realistically if in fact fascinatingly simulate the appearance of a `magic-wand` device with which children will particularly enjoy wandering about their neighborhoods holding on Halloween-nite(October -31), and on Independence-nite(July-4) when children often enjoy holding hotly burning `Sparklers`, which although exciting to youth, can be dangerous by inadvertantly setting afire their clothing.\nAccordingly, the invention as is now to be set forth, is presently being prepared for production by my San Diego, Calif. based `LightUp Products Co., Ltd*.` under the appropriate tradename of \"MagicWand\"-a type of safety-sparkler product."}
{"text": "1. Technical Field\nThe present disclosure relates to the application of buttress materials used in joining body tissue and attachment systems for attaching buttress materials to surgical instruments. More particularly, the present disclosure relates to systems and methods of temporarily attaching staple line buttress materials to the anvil and/or staple containing cartridge of a surgical stapling instrument.\n2. Background of Related Art\nSurgical stapling devices are employed by surgeons to sequentially or simultaneously apply one or more rows of fasteners, e.g., staples or two-part fasteners, to body tissue for the purpose of joining segments of body tissue together. Such devices generally consist of a pair of jaws or finger-like structures between which the body tissue to be joined is placed. When the stapling device is actuated, or “fired”, a firing member or members contact staple drive members in one of the jaws. The staple drive members push the surgical staples through the body tissue and into an anvil in an opposite jaw which crimps the staples closed. If tissue is to be removed or separated, a knife blade can be provided in the jaws of the device to cut the tissue between the lines of staples.\nWhen stapling relatively thin diseased, or fragile tissues, a buttress may be used to seal the staple line against air or fluid leakage. Additionally, a buttress may be used to reinforce the staple line against the tissue to prevent tears in the tissue or pulling of the staples through the tissue. A layer of buttress material is placed against the tissue and the tissue is stapled in conventional manner. It is known to position the buttress material on the stapling instrument itself prior to stapling the tissue. An exemplary example of this is disclosed in U.S. Pat. No. 5,542,594 to McKean et al. In McKean et al. the buttress material is slipped over the jaw of the stapler. The stapler is then actuated to staple the subject tissue and secure the buttress material between the tissue and staple line to reinforce the tissue and staple line.\nIt would be desirable to provide a system of attaching a limited amount of buttress material to a surgical instrument such that the material is localized on the jaws of the surgical instrument. It would also be desirable to provide a system of attaching buttress material to a surgical instrument that facilitates the release of the buttress material from the surgical instrument. Additionally, while the method disclosed in McKean is useful for open surgical procedures, it would be desirable to provide a system of attaching buttress material to the jaws of an endoscopic instrument. It would further be desirable to provide a system of attaching buttress material to the anvil and staple containing head of a circular surgical stapling instrument."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication apparatus.\n2. Description of the Related Art\nAn antenna which is used for radio communication radiates electromagnetic waves (radio waves) for radio communication after receiving input of a transmission signal (electric voltage or electric current) of a predetermined frequency. Conversely, the antenna also receives electromagnetic waves and converts into a reception signal (electric voltage or electric current). Here, when the antenna has a length corresponding to a wavelength which is used for the radio communication, receiving sensitivity and transmitting efficiency of electromagnetic waves can be improved.\nThus, the length of an antenna element (an antenna line) normally tends to be set at an approximate integral multiple of a quarter wavelength of the electromagnetic waves. In particular, the length of such an antenna element is an important factor for designing an antenna which is used for middle-range communication or long-range communication. It is similar to short-range communication as well. On the other hand, as disclosed in Japanese Patent Application Laid-Open No. 2007-36439, a technology of human body communication which utilizes a human body as a transmission path has been recently developed as for short-range communication. However, since a human body is utilized as the transmission path, the communication distance is limited. Therefore, it is difficult to apply the human body communication technology to middle-range communication or long-range communication."}
{"text": "1. Field of the Invention\nThis invention relates to a series of novel heterocyclic compounds that are useful in the treatment of hyperproliferative diseases, such as cancer and inflammation, in mammals. This invention also relates to a method of using such compounds in the treatment of hyperproliferative diseases in mammals, especially humans, and to pharmaceutical compositions containing such compounds.\n2. Description of the State of the Art\nCell signaling through growth factor receptors and protein kinases is an important regulator of cell growth, proliferation and differentiation. In normal cell growth, growth factors, through receptor activation (i.e. PDGF or EGF and others), activate MAP kinase pathways. One of the most important and most well understood MAP kinase pathways involved in normal and uncontrolled cell growth is the Ras/Raf kinase pathway. Active GTP-bound Ras results in the activation and indirect phosphorylation of Raf kinase. Raf then phosphorylates MEK1 and 2 on two serine residues (S218 and S222 for MEK1 and S222 and S226 for MEK2) (Ahn et al., Methods in Enzymology, 2001, 332, 417-431). Activated MEK then phosphorylates its only known substrates, the MAP kinases, ERK1 and 2. ERK phosphorylation by MEK occurs on Y204 and T202 for ERK1 and Y185 and T183 for ERK2 (Ahn et al., Methods in Enzymology, 2001, 332, 417-431). Phosphorylated ERK dimerizes and then translocates to the nucleus where it accumulates (Khokhlatchev et al., Cell, 1998, 93, 605-615). In the nucleus, ERK is involved in several important cellular functions, including but not limited to nuclear transport, signal transduction, DNA repair, nucleosome assembly and translocation, and mRNA processing and translation (Ahn et al., Molecular Cell, 2000, 6, 1343-1354). Overall, treatment of cells with growth factors leads to the activation of ERK1 and 2 which results in proliferation and, in some cases, differentiation (Lewis et al., Adv. Cancer Res., 1998, 74, 49-139).\nIn proliferative diseases, genetic mutations and/or overexpression of the growth factor receptors, downstream signaling proteins, or protein kinases involved in the ERK kinase pathway lead to uncontrolled cell proliferation and, eventually, tumor formation. For example, some cancers contain mutations which result in the continuous activation of this pathway due to continuous production of growth factors. Other mutations can lead to defects in the deactivation of the activated GTP-bound Ras complex, again resulting in activation of the MAP kinase pathway. Mutated, oncogenic forms of Ras are found in 50% of colon and >90% pancreatic cancers as well as many others types of cancers (Kohl et al., Science, 1993, 260, 1834-1837). Recently, bRaf mutations have been identified in more than 60% of malignant melanoma (Davies, H. et al., Nature, 2002, 417, 949-954). These mutations in bRaf result in a constitutively active MAP kinase cascade. Studies of primary tumor samples and cell lines have also shown constitutive or overactivation of the MAP kinase pathway in cancers of pancreas, colon, lung, ovary and kidney (Hoshino, R. et al., Oncogene, 1999, 18, 813-822). Hence, there is a strong correlation between cancers and an overactive MAP kinase pathway resulting from genetic mutations.\nAs constitutive or overactivation of MAP kinase cascade plays a pivotal role in cell proliferation and differentiation, inhibition of this pathway is believed to be beneficial in hyperproliferative diseases. MEK is a key player in this pathway as it is downstream of Ras and Raf. Additionally, it is an attractive therapeutic target because the only known substrates for MEK phosphorylation are the MAP kinases, ERK1 and 2. Inhibition of MEK has been shown to have potential therapeutic benefit in several studies. For example, small molecule MEK inhibitors have been shown to inhibit human tumor growth in nude mouse xenografts, (Sebolt-Leopold et al., Nature-Medicine, 1999, 5 (7), 810-816; Trachet et al., AACR Apr. 6-10, 2002, Poster #5426; Tecle, H., IBC 2nd International Conference of Protein Kinases, Sep. 9-10, 2002), block static allodynia in animals (WO 01/05390) and inhibit growth of acute myeloid leukemia cells (Milella et al., J. Clin. Invest., 2001, 108 (6), 851-859).\nSmall molecule inhibitors of MEK have been disclosed, including in U.S. Patent Publication Nos. 2003/0232869, 2004/0116710, and 2003/0216460, and U.S. patent application Ser. Nos. 10/654,580 and 10/929,295, each of which is hereby incorporated by reference. At least fifteen additional patent applications have appeared in the last several years. See, for example: U.S. Pat. No. 5,525,625; WO 98/43960; WO 99/01421; WO 99/01426; WO 00/41505; WO 00/42002; WO 00/42003; WO 00/41994; WO 00/42022; WO 00/42029; WO 00/68201; WO 01/68619; WO 02/06213; WO 03/077914; and WO 03/077855."}
{"text": "After devices or circuits are fabricated on a semiconductor wafer and similar substrates, the individual devices or circuits may be separated by scribing the wafer with a scribing tool, such as a sharply pointed scribe tip, and cleaving the wafer along the scribed lines.\nCurrent apparatus used to scribe a semiconductor wafer do not include a means or method of determining where the outer peripheral edge of the wafer is located before starting to scribe the wafer. Exemplary of prior art approaches are U.S. Pat. No. 3,094,785 and U.S. Pat. No. 5,820,006.\nIn these and similar apparatus, the scribe tip is placed outside the edge of the wafer and at a height just lower than the height of the wafer. The scribe tip is moved laterally relative to the wafer until it comes into contact with the wafer edge. It then continues across the wafer, forming a scribe line until it falls off the other side of the wafer.\nUsing this approach to wafer scribing causes two problems. First, when the scribe tip hits the leading edge of the wafer, damage to the wafer in the form of chips or cracks can occur. The second problem is that when the scribe tip hits the edge of the wafer it is knocked upward by the edge of the wafer and then oscillates up and down as the mechanism tries to dampen the mechanical perturbations in an attempt to maintain a constant scribe force."}
{"text": "Modern high-density integrated circuits (ICs) are known to be vulnerable to damage from the electrostatic discharge (ESD) from a charged body (human or otherwise) as the charged body physically contacts the IC. ESD damage occurs when the amount of charge exceeds the capability of the electrical conduction path through the IC. The typical ESD failure mechanisms include thermal runaway resulting in junction shorting, and dielectric breakdown resulting in gate-junction shorting in the metal-oxide-semiconductor (MOS) context.\nAn IC may be subjected to a damaging ESD event in the manufacturing process, during assembly, testing, or during runtime of the system in which the IC is installed. Some ESD protection schemes use active clamp circuits to shunt ESD current between an input/output (I/O) and ground pad and thereby protect internal IC element nodes that are connected to bond pads from ESD damage. An active clamp circuit may include a Metal Oxide Semiconductor Field Effect Transistor (MOSFET) coupled between the I/O and ground pad. This MOSFET is relatively large and often occupies 60% or more of total area of the ESD protection circuit. Each of the IC's pins generally includes its own ESD protection circuit with its own large MOSFET."}
{"text": "The present invention relates to military training equipment, and more particularly, to a system for automatically aligning a laser transmitter on a rifle for subsequent use by a soldier in war games.\nFor many years the armed services of the United States have trained soldiers with a multiple integrated laser engagement system (MILES). A laser small arms transmitter (SAT) is affixed to the stock of a rifle such as an M16. Each soldier carries detectors on his helmet and on a body harness adapted to detect a laser \"bullet\" hit. The soldier pulls the trigger of his or her rifle to fire a blank to simulate the firing of an actual round and an audio sensor triggers the SAT.\nIt is necessary to align the SAT so that the soldier can accurately hit the target once he or she has it located in the conventional rifle sights. In the past an early version of the SAT was bolted to the rifle stock and the mechanical sights of the weapon were adjusted to align with the laser beam. The disadvantage of this approach is that the mechanical weapon sights must be readjusted in order to use the rifle with live rounds. To overcome this disadvantage the conventional SAT now in use incorporates mechanical linkages for changing the orientation of the laser.\nThe prior art small arms alignment fixture (SAAF) used by the U.S. Army for alignment of the conventional MILES SAT consists of a complex array of one hundred forty-four detectors which are used in conjunction with thirty-five printed circuit boards to determine where the laser hits with respect to a target reticle. The difficulty in using the prior art SAAF is that the soldier aims his or her weapon at the array which is twenty-five meters away without the use of a stable platform. In many cases, the soldier fires his or her weapon in a manner which results in the aim point not being at the desired location. The fact that the array is located twenty-five meters away from the soldier introduces visibility limitations due to snow, fog, wind and poor lighting conditions at sunrise or dusk. The prior art SAAF calculates the number of error \"clicks\" in both azimuth and elevation. The number of clicks is then displayed on the prior art SAAF using four sets of electro-mechanical display indicators. The soldier must then turn his conventional SAT's adjustors the corresponding number of clicks in the correct direction. He or she must then aim and fire the weapon again and make additional corresponding adjustments. This iterative process continues until the soldier obtains a zero indication on the prior art SAAF. This is a very time consuming and tedious process due to normal aiming errors incurred each time the soldier has to reacquire the target reticle. It is not uncommon for a soldier to take fifteen minutes to align his or her weapon to the best of his or her ability and still not have it accurately aligned.\nNot only is the alignment process utilizing the prior art SAAF time consuming, it is also expensive because a large amount of blank ammunition must be used. The laser of a conventional SAT will not fire without a blank cartridge being ignited or by using a special dry fire trigger cable. The prior art SAAF does not support optical sights, different small arms weapon types, nor night vision devices. Nor does the prior art SAAF accurately verify the laser beam energy and encoding of the received laser beam.\nIt would therefore be desirable to provide an improved alignment system for a small arms SAT which would eliminate the need to utilize a large target array. Such a system would also preferably automatically adjust the SAT for more rapid and accurate alignment. In addition, preferably such a system would require only a single target sighting and would accommodate different small arms such as automatic weapons, sniper rifles, and so forth. Not only do these small arms have different gun stocks, but in addition, the laser output of their SATs have different powers and codings to enable the manworn portion of a MILES system to discriminate between hits made by different small arms."}
{"text": "Cytotoxic T lymphocytes (CTL) are important effectors of immunity to pathological conditions such as microbial infections and cancer in that they have the ability to destroy “abnormal” cells such as those infected by intracellular pathogens or tumor cells. CTL recognize these abnormal cells as “foreign” via cell-surface antigen-specific T cell receptors (TCR), which bind to antigenic peptide epitopes complexed with major histocompatibility complex (MHC) molecules (most commonly MHC class I molecules) on the surface of the abnormal cells. CTL peptide epitopes are generally derived from intracellular proteins such as microbial (e.g., viral) proteins and tumor-associated antigens that are: (i) synthesized within the cell; (ii) proteolytically processed within the cell; and (iii) transported by the TAP transporter molecules into the endoplasmic reticulum (ER) [Germain et al. (1994) Cell 76:287; York et al. (1996) Annual Review of Immunology 14:369; Heemels et al. (1995) Annu Rev Biochem 64:463; Momburg et al. (1998) Adv Immunol 68:191]. Once in the ER, the CTL peptide epitopes can be further “trimmed” by the action of peptidases into an optimal size for association with MHC class I molecules; after complexing of CTL peptide epitopes with an appropriate MHC class I molecule, the resulting complex is exported to the cell-surface where it is available for recognition by a relevant TCR [Snyder et al. (1994) J Exp Med 180:2389; Roelse et al. (1994) J Exp Med 180:1591; Mo et al. (1999) J Immunol 163:5851; Paz et al. (1999) Immunity 11:241; Powis et al. (1996) Immunity 4:159; Yewdell et al. (2001) Curr Opin Immunol 13:13]. This process is known as the “MHC class I antigen-processing pathway.”\nThe requirement for intracellular processing of antigens creates some constraints on the ability to induce CTL responses using exogenous, non-infectious vaccines such as killed pathogens or protein antigens. These constraints arise because cells, e.g., antigen presenting cells (APC), are generally not very efficient at processing via the MHC class I antigen-processing pathway antigens that are not synthesized within the cell but are imported into the cell by any of the mechanisms by which such importation occurs."}
{"text": "Tang and coworkers first reported on high-performance organic light-emitting devices (OLEDs) in 1987 (Tang, C. W.; et al. Appl. Phys. Lett. 51, 913 (1987)). Their discovery was based on employing a multilayer structure containing an emitting layer and a hole transport layer of a suitable organic substrate. Alq3(q =deprotonated 8-hydroxyquinolinyl) was chosen as the emitting material and proven to be of high performance because (1) it can form uniform thin films under 1000 Å using vacuum deposition, (2) it is a good charge carrier and (3) it exhibits strong fluorescence. Since then, there has been a flourish of research on OLEDs and materials used in these devices. Indeed, nearly every large chemical company in the world with optoelectronic interests has demonstrated some level of interest in OLEDs. Clearly, OLED technology is heading directly and rapidly into the marketplace, as suggested in a commercial report by Stanford Resources (by David E. Mentley, “The Market Potential for Organic Light-Emitting Diode Displays,” Commercial Report. The attractiveness of OLEDs as it challenges traditional technologies such as cathode ray tubes (CRTs), liquid crystal displays (LCDs) and plasma displays is based on many features and advantages, including:                Low operating voltage,        Thin, monolithic structure,        Emits, rather than modulates light,        Good luminous efficiency,        Full color potential, and        High contrast and resolution.        \nOLED is a device built with organic semiconductors from which visible light can be emitted upon electrical stimulation. The basic heterostructure of an OLED is described in FIG. 1.\nThe layers may be formed by evaporation, spin-casting or chemical self-assembly. The thickness ranges from a few monolayers (self-assembled films) to about 1000 to 2000 Å. Such devices whose structure is based on the use of layers of organic optoelectronic materials generally rely on a common mechanism leading to optical emission, namely, the radiative recombination of a trapped charge. Under a DC bias, electrons are injected from a cathode (usually Ca, Al, Mg—Ag) and holes are injected from an anode (usually transparent indium tin oxide (ITO)) into the organic materials, where they travel in the applied field across the electron transporting layer (ETL) and the hole transporting layer (HTL) respectively until they meet, preferably on molecules in the emitting layer, and form a luminescent excited state (Frenkel exciton) which, under certain conditions, experiences radiative decay to give visible light. The electroluminescent material may be present in a separate emitting layer between the ETL and the HTL in what is referred as a multi-layer heterostructure. In some cases, buffer layers and/or other functional layers are also incorporated to improve the performance of the device. Alternatively, those OLEDs in which the electroluminescent emitters are the same materials that function either as the ETL or HTL are referred to as single-layer heterostructures.\nIn addition to emissive materials that are present as the predominant component in the charge carrier layers (HTL or ETL), other efficient luminescent material(s) may be present in relatively low concentrations as a dopant in these layers to realize color tuning and efficiency improvement. Whenever a dopant is present, the predominant material in the charge carrier layer may be referred to as a host. Ideally, materials that are present as hosts and dopant are matched so as to have a high level of energy transfer from the host to the dopant, and to yield emission with a relatively narrow band centered near selected spectral region with high-efficiency and high-brightness.\nWhile fluorescent emitters with high luminescence efficiencies have been extensively applied as dopant in OLEDs, phosphorescent emitters have been neglected in this domain. However, the quantum efficiency of an electrofluorescence device is limited by the low theoretical ratio of singlet exciton (25%) compared to triplet exciton (75%) upon electron-hole recombination from electrical excitation. In contrast, when phosphorescent emitters are employed, the potentially high energy/electron transfer from the hosts to the phosphorescent emitters may result in significantly superior electroluminescent efficiency (Baldo, M. A.; et al. Nature 395, 151 (1998) and Ma, Y. G.; et al. Synth. Met. 94, 245 (1998)). Several phosphorescent OLED systems have been fabricated and have indeed proven to be of relative high-efficiency and high-brightness.\nIt is desirable for OLEDs to be fabricated using materials that provide electrophosphorescent emission corresponding to one of the three primary colors, i.e., red, green and blue so that they may be used as a component layer in full-color display devices. It is also desirable that such materials are capable of being deposited as thin films using vacuum deposition techniques, which has been prove to be a common method for high-performance OLED fabrication, so that the thickness of the emitting layer can be precisely controlled.\nPresently, the highest efficiencies and brightness have been obtained with green electrophosphorescent devices (15.4±0.2% for external quantum efficiency and almost 100% for internal efficiency, 105 Cd/m2 for maximum luminance) using Ir(ppy)3 (ppy=deprotonated 2-phenylpyridine) as emitter (Adachi, C.; et al. Appl Phys. Lett. 77, 904 (2000)). An OLED emitting saturated red light based on the electrophosphorescent dopant Pt(OEP) (H2OEP=octaethylporphyrin) has also been published and patented (Burrows, P.; et al. U.S. Pat. No. 6,048,630) but the maximum luminance is only around 500 Cd m−2. A relevant patent is the use of the cyclometalated platinum(II) complex Pt(thpy)2 (thpy=deprotonated 2-(2-thioenyl)pyridine) as dopant and PVK (poly(N-vinyl)carbazole) as host in a orange OLED (Lamansky, S.; et al. WO Pat. No. 00/57676). However, the Pt(ID) complex used by the inventors was not stable for sublimation or vacuum deposition, thus a spin-casting method was applied, which led to higher driving voltages, quantum efficiency of 0.11% and luminance of 100 Cd/m2 were obtained at 22 V."}
{"text": "Our unique rotary mineral breaker is described and claimed in our U.S. patent specification No. 3,970,257. While this apparatus is extremely satisfactory in use a problem can arise where one of the hardened wear tips overwhich the minerals pass fails and the apparatus is still maintained in use. In such circumstances a section of the rotor can be worn away comparatively quickly leading to a significant maintenance operation."}
{"text": "1) Field of the Invention\nThe present invention relates to a system for detecting a position of a device such as a capsule endoscope, which is swallowed and passes naturally through the digestive tract of the subject, by a position transducer that is disposed on the outside of the subject and obtains information of the position of the device in the subject.\n2) Description of the Related Art\nIn recent years, in the field of endoscopes, a swallowable capsule endoscope has been proposed. The capsule endoscope has an image capturing function and a radio communication function. The capsule endoscope has the function of traveling in the body cavity, for example, in the organs such as the stomach and the small intestine with peristalsis of the organs and sequentially capturing images for a period of time since the capsule endoscope is swallowed from the mouth of a subject for inspection (examination) until it is naturally excreted.\nImage data captured in the body by the capsule endoscope as the capsule endoscope travels in the body cavity is sequentially transmitted by radio communication to the outside and stored into a memory provided on the outside. The subject can freely move throughout the period after he/she swallows the capsule endoscope until it is excreted by carrying a receiver having a radio communication function and a storing function. After the capsule endoscope is excreted, a doctor or nurse can display the images of the organs on a display based on the image data stored in the memory and make a check.\nA capsule endoscope has been proposed in which the receiver has the function of detecting the position of the capsule endoscope in the subject to capture, for example, an endoscope image of a specific organ in the subject. As an example of a capsule endoscope system having the position detecting function, a capsule endoscope system using the radio communication function provided in the capsule endoscope is known. Specifically, the system has a configuration that a receiver provided on the outside of a subject has a plurality of antenna elements, and has the function of receiving a radio signal transmitted from the capsule endoscope by the plurality of antenna elements and, based on intensities received by the antenna elements, detecting the position of the capsule endoscope in the subject (see Japanese Patent Application Laid-open No. 2003-19111, for example)."}
{"text": "Electrode catheters have been in common use in medical practice for many years. They are used to stimulate and map electrical activity in the heart and to ablate sites of aberrant electrical activity.\nIn use, the electrode catheter is inserted into a major vein or artery, e.g., femoral vein, and then guided into the chamber of the heart which is of concern. In some medical procedures, energy is imparted to body tissue locally, in a concentrated dose, and it is desirable to cool the treatment area in order to reduce collateral tissue damage. For example, cardiac ablation therapy is used to treat arrhythmias by heating tissue with radio-frequency (RF) electrical energy to create non-conducting lesions in the myocardium. It has been found that cooling the area of the ablation site reduces tissue charring and thrombus formation. Catheters with irrigated distal tips are known as part of integrated ablation system. Typically, a metal catheter tip, which is energized with RF current to ablate the tissue, has a number of irrigation holes, distributed circumferentially around the tip, for irrigation of the treatment site. A pump coupled to the catheter delivers saline solution to the catheter tip, and the solution flows out through the holes during the procedure in order to cool the catheter tip and the tissue.\nIn certain regions of the heart, for example, in the ventricles where tissue is thicker, the creation of transmural lesions can be challenging. Deep lesions typically require higher RF energy but higher RF energy can lead to undesirable steam pops. Thus, there is a desire to create deeper lesions by increasing electrode/tissue contact area but without increasing the size of the catheter itself.\nCatheters with flexible tips are known. U.S. Pat. No. 5,720,719 describes a catheter having a probe end that includes a malleable tube and a flexible tube. U.S. Patent Publication No. 2014/0121657, whose disclosure is incorporated herein by reference, describes a medical probe having a deformable distal end that includes a flexible and porous material. The flexible and porous material may include a conductive material. An electrical conductor can be coupled to the flexible and porous material so as to convey RF energy to the deformable distal end, and the RF energy can be conveyed to tissue by the deformable distal end conveying the RF energy to the tissue. The medical probe may include means for inflating the deformable end which may include conveying a fluid that irrigates the tissue through pores of the deformable distal end. The means for inflating the deformable distal end may include conveying the fluid the fluid so as to generate a mechanical force sufficient to inflate the deformable distal end. A contact area between the deformable distal and the tissue can increase upon pressing the deformable distal end against the tissue.\nU.S. Pat. No. 8,249,685 is directed to an apparatus for mapping and/or ablating tissue that includes a braided conductive member that may be inverted to provide a ring shaped surface. When a distal tip of the braided conductive member is retracted within the braided conducive member, the lack of protrusion allows the ring-shaped surface to contact a tissue wall such as a cardiac wall. In an undeployed configuration, the braided conductive member is longitudinally extended, and in a deployed configuration, the distal end of the braided conductive member is retracted to invert the braided conductive member.\nThe descriptive above is presented as a general overview of related art in this field and should be not be construed as an admission that any of the information it contains constitutes prior art against the present patent application."}
{"text": "1. Field of the Invention\nThe present invention relates to disk drives for computer systems. More particularly, the present invention relates to servo writing a disk drive from spiral tracks by generating a time-stamped sync mark detect signal processed by timing recovery firmware.\n2. Description of the Prior Art\nWhen manufacturing a disk drive, product servo sectors 20–27 are written to a disk 4 which define a plurality of radially-spaced, concentric data tracks 6 as shown in the prior art disk format of FIG. 1. Each product servo sector (e.g., servo sector 24) comprises a preamble 8 for synchronizing gain control and timing recovery, a sync mark 10 for synchronizing to a data field 12 comprising coarse head positioning information such as a track number, and servo bursts 14 which provide fine head positioning information. During normal operation the servo bursts 14 are processed by the disk drive in order to maintain a head over a centerline of a target track while writing or reading data. In the past, external servo writers have been used to write the product servo sectors 20–27 to the disk surface during manufacturing. External servo writers employ extremely accurate head positioning mechanics, such as a laser interferometer, to ensure the product servo sectors 20–27 are written at the proper radial location from the outer diameter of the disk to the inner diameter of the disk. However, external servo writers are expensive and require a clean room environment so that a head positioning pin can be inserted into the head disk assembly (HDA) without contaminating the disk. Thus, external servo writers have become an expensive bottleneck in the disk drive manufacturing process.\nThe prior art has suggested various “self-servo” writing methods wherein the internal electronics of the disk drive are used to write the product servo sectors independent of an external servo writer. For example, U.S. Pat. No. 5,668,679 teaches a disk drive which performs a self-servo writing operation by writing a plurality of spiral tracks to the disk which are then processed to write the product servo sectors along a circular path. Each spiral track is written to the disk as a high frequency signal (with missing bits), wherein the position error signal (PES) for tracking is generated relative to time shifts in the detected location of the spiral tracks. In addition, the '679 patent generates a servo write clock by synchronizing a phase-locked loop (PLL) to the missing bits in the spiral tracks. However, in the '679 patent the timing recovery measurement for updating the PLL is generated in hardware using complex and expensive timing circuitry for measuring the phase difference between the detected missing bits and the servo write clock.\nThere is, therefore, a need for a simple, less expensive technique for synchronizing a servo write clock when servo writing a disk drive by servoing on spiral tracks."}
{"text": "The Android system is a new generation of smart mobile terminal system. In more and more mobile terminals, the Android system is used as an operation system, and the users of the mobile terminals adopting the Android system also increasingly grow. The mobile terminal with the Android system provides strong network functions, and the users can Access Internet very conveniently.\nThere are two basic concepts in the Android system need to be introduced firstly, one is APN, and the other is PDP. Wherein:\nAccess Point Name (APN): the logical name of one Gateway GPRS Support Node (GGSN) and the external network; corresponding APN configuration is required when a PDP is activated.\nPacket Data Protocol (PDP): an interaction channel between a mobile terminal and the network; one PDP enables the mobile terminal to own an IP address, and it is corresponding to one network interface in the mobile terminal.\nThrough studies the inventor finds that currently, the mobile terminal adopting the Android system does not support usage of multiple APNs simultaneously while performing the network related services. That is to say, at one timing, in the mobile terminal adopting the Android system there are only one APN in the being used state, and correspondingly, there is just one PDP exists in the mobile terminal adopting the Android system. In case that the APN existed in the mobile terminal with the Android system is an APN for a dedicated service, e.g. an APN specialized for the Internet service, and corresponding PDP can only process the Internet service. When the mobile terminal being processing the Internet service receives a request for processing the Multimedia Messaging Service (MMS), a processing method on the current Android system firstly deactivates a PDP presently in the activated state and being processing the Internet service, then activates a PDP capable of processing the MMS according to the configuration of the APN capable of processing the MMS, and recovers the former PDP for processing the Internet service after the processing of the MMS is ended. It is clear that in the current Android system, when a PDP corresponding to a dedicated APN is processing a service, if a request for processing a different type of service is received, the processing of the current service shall be interrupted to process the new service, which greatly influences the processing of the original service, reduces the processing speed of the mobile terminal, and decreases the user experience."}
{"text": "Electronic packages can be molded in a plastic compound to improve the interconnect reliability between a semiconductor die and a substrate and to help keep the substrate flat for assembly onto a circuit board. The molding compound and the substrate can have different coefficients of thermal expansion, glass transition temperatures, shrink rates, and/or mechanical rigidity, which can cause warping and/or stresses on the semiconductor die and/or on the substrate when there are room temperature changes or during temperature cycling.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present disclosure as set forth in the remainder of the present application with reference to the drawings."}
{"text": "For many years, ground fault circuit interrupters (GFCIs) have served an important function in the electrical power distribution system of private homes, businesses, and other structures in preventing users from experiencing an electrical shock under many circumstances. Typically, a GFCI is implemented within an electrical power receptacle and performs its function by disconnecting power at a receptacle connection pair of the receptacle in response to detecting small differences in current at the receptacle carried over a live or “hot” current path and a neutral current path that may be indicative of a “ground fault,” in which at least some current is inadvertently being directed to ground from the live or neutral current path. Typically, such detection is performed using a current transformer through which the live and neutral paths are routed within the GFCI. In addition, the GFCI may also provide power to one or more additional receptacles on the same electrical power circuit.\nTo perform its function properly, a GFCI receptacle typically provides a line connection pair for live and neutral wires received from a circuit breaker of an electrical panel, and also provides a load connection pair for live and neutral wires for one or more additional receptacles of the same electrical power circuit, in addition to the receptacle connection pair, so that the GFCI may disconnect both the receptacle connection pair and the load connection pair from the line connection pair.\nGenerally speaking, miswiring of the GFCI, in which the live and neutral wires from the circuit breaker panel are mistakenly coupled to the live and neutral lines of the load connection pair, and the live and neutral wires connected to the other power receptacles are connected to the line connection pair, is an all-too-common occurrence. In cases of such miswiring, most GFCIs are unable to detect a ground fault at the GFCI receptacle connection pair since the current transformer is not positioned between the load connection pair and the receptacle connection pair, and thus cannot detect the resulting difference in current in the live and neutral current paths at that point.\nIn addition, other electrical safety devices, such as, for example, arc fault circuit interrupters (AFCIs), may also be susceptible to a similar miswiring problem."}
{"text": "The present invention relates to an improved needle-less parenteral fluid injector which is more efficient and is safer for both the health-care worker as well as the patient."}
{"text": "Multimedia messaging systems are commercially available for such applications as voice mail and short message services over fixed telephone, radio telephone, paging and other types of communications networks. One such commercially available multimedia messaging system is the so-called MXE system manufactured by the present assignee. Throughout this specification, where reference is made to Ericsson's MXE system, that term shall mean either the commercially available Ericsson MXE system or such other multimedia messaging systems as may be available.\nTraditionally, multimedia messaging systems are available with voice mail, fax mail, e-mail, notification services, short message services (alphanumeric messages displayed on mobile phones), and other customized services for the subscriber. The multimedia messaging system is the hardware and software that provides the infrastructure for these messaging applications across all types of telecommunications networks, such as fixed, cellular, private, paging, personal communication services (PCS), etc.\nMultimedia messaging systems communicate with the network via multimedia interfaces (MMIs). An example multimedia interface (MMI) supporting present messaging applications for multimedia messaging systems is disclosed, for example, in Hulen et al., U.S. Pat. No. 5,497,373, which is incorporated herein by reference.\nThe MXE, like other multimedia messaging systems, supports a limited number of subscribers. This is due to at least three inherent bottlenecks in the system, namely the DMA transfers (i.e., the playing of voice prompts and the recording of voice messages), the ability to keep up with a remote telephony subsystem, and the disk access (such as subscriber lookups and message storage/retrieval). These bottle-necks are in fact the normal types of bottlenecks that you would expect to see on a system of this type, for example, on a system built around a single off-the-shelf UNIX workstation.\nThe multimedia messaging system, from the subscribers' view, provides personalized communications available in a variety of different formats to permit the subscriber to access messages in a variety of different ways. For example, when subscriber may otherwise be out of reach of traditional telephone systems, using the multimedia messaging system, others can leave voice messages, send faxes, or send e-mail to the subscriber via the multimedia messaging system which automatically notifies the subscriber that messages are waiting. When it is convenient, the subscribers then call a single number to retrieve all of the messages.\nFrom the network perspective, the multimedia messaging system provides integrated messaging services to market to the subscribers, provides a single platform for developing additional services to provide to the subscribers, and enhances revenue streams by increasing call completion ratios.\nDemand for multimedia messaging systems services is increasing, causing the bottlenecks inherent in multimedia messaging systems to become increasingly significant. It therefore becomes necessary to add additional multimedia messaging systems to existing ones. This is traditionally done by adding systems containing sophisticated signaling interfaces to the telecom environment. One such example system is the so-called \"LNX\" manufactured by Excel, Inc. of Massachusetts. The LNX is a physical box having ports and a cage of circuit boards providing the interfacing to the telecom environment through an additional controller software module (sometimes referred to as an \"LNX controller\"). In this disclosure, the terms \"LNX\" and \"LNX Controller\" should be interpreted to broadly refer to sophisticated interfacing equipment whether marketed under the \"LNX\" name or not. In any event, expanding system capacity by adding new LNX boxes can be rather expensive and inefficient.\nIt thus becomes important to be able to add capacity to an existing multimedia messaging system, without the expense of additional telephony-equipped workstations."}
{"text": "1. Field of the invention\nThe present invention relates to the field of multilevel integrated circuits, and more specifically, to the field of ULSI multilevel interconnection structures and their methods of fabrication.\n2. Discussion Of Related Art\nModem integrated circuits are made up of literally millions of active devices such as transistors and capacitors. These devices are initially isolated from one another but are later interconnected together to formed functional circuits. The quality of the interconnection structure drastically affects the performance and reliability of the fabricated circuit. Interconnections are increasingly determining the limits of performance and density of modern ultra large scale integrated (ULSI) circuits.\nFIG. 1 is a cross-sectional illustration of an interconnection structure which is presently used in the semiconductor industry. Formed in a silicon substrate or well 102 are active devices (not shown) such as transistors and capacitors. Interconnection lines 104 and 106, which are typically aluminum or aluminum alloys, are used to couple active devices into functional circuits. Metal interconnections 104 and 106 and substrate 102 are electrically isolated from one another by interlevel dielectric\"\"s (ILDs) 108 and 110, respectively. Electrical connections are made between metal interconnections 104 and 106 through the use of metal vias 112.\nThe interconnection structure of FIG. 1 experiences several problems. As integrated circuit dimensions decrease, in order to increase circuit density and complexity, vias are becoming smaller by the square of the dimension decrease. Such small vias can cause both reliability and performance problems in a fabricated integrated circuit. Reliability problems are caused by high concentrations of current or current crowding at the corners 114 of via 112, as shown in FIG. 1. Current crowding in area 114 can cause self-heating effects which can cause electromigration. Electromigration can cause voids and open circuits and other reliability problems. Performance of small dimensioned vias is decreased due to an increase in contact resistance caused by a reduction of the interfacial contact area between via and interconnections. It is to be appreciated that contact resistance is becoming the major portion of the total interconnection resistance.\nAnother problem with the interconnection structure of FIG. 1 is the aluminum metal used for interconnections 104 and 106. Aluminum interconnections are known to suffer from electromigration. Electromigration in interconnections are known to cause voids which can cause open circuits and failures. Aluminum interconnections are also known to form hillocks which can result in nonplanar interconnection structures and can cause interlevel shorts with interconnections formed above and intralevel shorts with neighboring interconnection lines from the same level of metalization. These problems are magnified as line widths decrease and packing densities increase as in the case of future ULSI circuits.\nThus, what is desired is a ULSI compatible interconnection structure which exhibits excellent performance and reliability.\nA novel, high performance, high reliability, high density, ULSI manufacturable, multilevel interconnection structure is described. The interconnection structure of the present invention is formed on a first insulating layer of an integrated circuit. A first multilayer interconnection comprising a first aluminum layer, a first refractory metal layer, and a second aluminum layer is formed on the first insulating layer. A second insulating layer is formed over the first multilayer interconnection. A conductive via of, for example, tungsten, aluminum, copper or TiN, is formed through the second insulating layer and into the multilayer interconnection wherein the via is formed down to at least the first refractory metal layer of the multilayer interconnection to form a high surface area contact with the first multilayer interconnection. A portion of the conductive via extends above the second insulating layer so that a high surface area contact can be made to a second interconnection which is formed on the second insulating layer and on and around the portion of the via extending above the second insulating layer.\nA goal of the present invention is to form an interconnection structure which exhibits excellent performance and reliability even when formed to the small dimensions necessary for modern ULSI circuits.\nAnother goal of the present invention is to provide an interconnection structure which provides decreased current crowding around vias.\nStill another goal of the present invention is to provide an interconnection structure which exhibits redundant current carrying capabilities.\nStill yet another goal of the present invention is to provide an interconnection structure which is resistant to electromigration and related failures.\nStill other goals and advantages of the present invention will become obvious from the detailed description of the present invention which follows."}
{"text": "1. Field of the Invention\nThe invention is concerned with polymers in aqueous colloidal dispersion, i.e., latexes, which are stabilized by isothiouronium cations.\n2. Description of the Prior Art\nU.S. Pat. No. 3,123,636 describes thiuronium halides which are viscous oils that are soluble in water and are useful as insect attractants and bactericides. The patent states that the compounds when shaken in water cause foaming, thus illustrating surfactant properties.\nU.S. Pat. No. 3,594,355 describes water-soluble isothiuronium salts of epihalohydrin polymers and copolymers. These water soluble salts are disclosed as highly useful as shrinkproofing agents for wool, as wet and dry strength paper additives, as flocculating agents, as cross-linkable coatings for textiles and as semipermeable membranes. There is no suggestion that any kind of latex can be prepared with the polymers.\nU.S. Pat. No. 3,642,879 describes S-(vinylbenzyl)isothiouronium compounds in which the isothiouronium moiety has the formula ##STR1## These compounds are useful as pesticides for the control of various plant, insects, and fungal pests and are monomers which can be polymerized. A latex of a copolymer of one of these compounds with styrene and butadiene is described. Such latexes, however, are useful and stable only in acidic media over a narrow pH range."}
{"text": "Many consumer electronic appliances, such as television set top boxes and remote control devices, are implemented as electronic systems incorporating integrated circuits. These electronic systems typically include integrated circuits forming the core digital logic circuitry and integrated circuits forming the input/output (I/O) interface circuitry. Conventional electronic systems employ multiple DC-to-DC converters to convert a main bus voltage from a power source supplying the system to one or more voltages necessary for driving these integrated circuits.\nSwitch mode power supplies or switching regulators, also referred to as DC to DC converters, are used to convert an input supply voltage to a desired output voltage at a voltage level appropriate for integrated circuits in an electronic system. For example, a 12 volts supply voltage provided to an electronic system may need to be reduced to 5 volts for supplying the I/O interface circuits and reduced to 1V for supplying the core digital logic circuits, especially if the core digital logic circuits are built using deep sub-micron integrated circuits. A switching regulator provides power supply function through low loss components such as capacitors, inductors, and transformers, and power switches that are turned on and off to transfer energy from the input to the output in discrete packets. A feedback control circuit is used to regulate the energy transfer to maintain a constant output voltage within the desired load limits of the circuit.\nGovernment agencies have started to implement standby power regulations, requiring that appliances with standby power function use no more than 1 watt of power in their standby power consuming mode. These “green” regulations require electronic appliances to have high power efficiency with low standby power consumption. Power efficiency is a measure of the fraction of power consumed which achieves functionality, rather than waste heat.\nWhen switching regulators are incorporated in electronic systems, standby power regulations demand high efficiency from the switching regulators when in standby mode, such as greater than 80% when drawing 10 mA at 5V during standby. Switching regulators are designed to deliver peak efficiency when driving some specific “normal” load. Switching regulators lose efficiency when operated in a wide range of load conditions, and the most drastic degradation occurs at light load conditions. Thus, it is often difficult for switching regulators to meet the standby power requirements in light load conditions.\nIn an electronic system including a digital core circuit and an I/O interface circuit, the digital core circuit typically demands high current but operates at a low power supply voltage (e.g. 1V) while the I/O interface circuit typically requires lower current but operates at a higher power supply voltage (e.g. 5V or 3.3V).\nThe digital core circuit requires a power supply capable of fast load transient response to handle varying load current transitions quickly. Since the digital core is typically turned off during standby mode, the efficiency of the digital core power supply is not critical for the standby regulations. Therefore, a DC-to-DC converter capable of fast load transient is typically selected, regardless of the converter's power efficiency at light load. In most cases, a pulse width modulation (PWM) switching regulator is used to supply the digital core circuit.\nOn the other hand, the I/O interface circuit remains turned on during the standby mode of operation. Thus, the I/O interface circuit requires a power supply capable of delivering high efficiency to a normal load as well as to a light load. In order to meet the requirements of the “green” regulations, a power supply for the I/O interface circuit needs to have high efficiency at the light load condition when the system is in standby mode. In most cases, a pulse frequency modulation (PFM) switch regulator is desired for high efficiency light load operation.\nAccordingly, to optimize both the performance requirement and the power efficiency demand of an electronic system, an electronic system may end up being designed to use different types of switching regulators for the digital core circuit and the I/O interface circuit. A manufacturer of an electronic system thus have to procure multiple DC-to-DC converter solutions, rendering the design and manufacturing process more complex and less cost effective."}
{"text": "Image encoding techniques have been applied to many video devices that are familiarly used. Target products to which the techniques are applied may be, for example, DVD for enjoying video contents, such as movies, hard disk recorders in which video contents, such as TV broadcasts, are to be recorded, digital TV encoding systems, DVD cameras, and mobile phones that can carry out image processing. As such, the image encoding techniques are applied very widely.\nNon Patent Literature 1 discloses (i) an image encoding device that encodes an image by use of a spatial correlation or a temporal correlation in the image, and (ii) an image decoding device that decodes encoded data that is encoded by the image encoding device. The following explains about techniques of these devices that use the spatial correlation, with reference to FIGS. 65 through 70.\n(Arrangement of Image Encoding Device 300)\nWith reference to FIG. 65, the following describes an image encoding device 300 according to a conventional technique.\nFIG. 65 is a block diagram showing essential part of an arrangement of the image encoding device 300 according to the conventional technique. The image encoding device 300 includes a difference computing section 1, an orthogonal transformation section 2, a quantization section 3, an entropy encoding section 304, an inverse quantization section 6, an inverse orthogonal transformation section 7, an addition computing section 8, a memory 9, an intra prediction section 310, and a prediction mode determining section 311 (FIG. 65). The following describes each of the constituent components of the image encoding device 300.\n(Difference Computing Section 1)\nThe difference computing section 1 outputs prediction residual data, which is a difference between a target block to be encoded (M×M pixel block) and a predicted image generated from the intra prediction section 310.\n(Orthogonal Transformation Section 2 and Inverse Orthogonal Transformation Section 7)\nThe orthogonal transformation section 2 carries out orthogonal transformation with respect to the prediction residual data received from the difference computing section 1, and then outputs the resultant. The inverse orthogonal transformation section 7 carries out inverse orthogonal transformation with respect to an orthogonal transform coefficient of the prediction residual data received from the inverse quantization section 6, and outputs the resultant. As a method of the orthogonal transformation and the inverse orthogonal transformation, discrete cosine transform, Hadamard transform, discrete Fourier transform, discrete sine transform, Haar transform, slant transform, or Karhunen-Loeve transform can be used.\n(Quantization Section 3 and Inverse Quantization Section 6)\nThe quantization section 3 carries out quantization with respect to an orthogonal transformation coefficient of the prediction residual data received from the orthogonal transformation section 2, and then outputs the resultant. The inverse quantization section 6 carries out inverse quantization with respect to a quantization coefficient of the prediction residual data received from the quantization section 3, and then outputs the resultant. The quantization and the inverse quantization are carried out by use of scalar quantization or vector quantization.\n(Entropy Encoding Section 304)\nThe entropy encoding section 304 carries out entropy encoding with respect to the prediction residual data and information on an encoding mode, such as a block type, prediction mode information, and a quantization parameter. The “entropy encoding” indicates: (a) variable length coding, such as arithmetic coding, Huffman coding, and Golomb coding; or (b) fixed length coding.\n(Addition Computing Section 8)\nThe addition computing section 8 adds the prediction residual data subjected to the inverse quantization and the inverse orthogonal transformation, to a predicted image generated by the intra prediction section 310 so as to generate a locally decoded image. The addition computing section 8 then output the locally decoded image.\n(Memory 9)\nThe memory 9 receives and then stores the locally decoded image therein.\n(Intra Prediction Section 310)\nThe intra prediction section 310 carries out intra prediction indicated by prediction mode information, by use of a locally decoded image stored in the memory 9, and generates a predicted image. The “intra prediction” indicates prediction in a screen or prediction in a frame, that is, prediction carried out by use of a spatial correlation in the image. At this time, the intra prediction is carried out on a target block (M×M pixel block) to be encoded, by use of 9 types of predetermined prediction methods respectively indicated by Modes 0 Through 8. FIG. 66 shows examples of prediction modes. FIG. 66 shows prediction modes for use in the intra prediction according to the conventional technique. In the examples in FIG. 66, a block size is a 4×4 pixel block. Further, pixels A through M are pixels which have been already encoded and which are used for prediction of the target block to be encoded.\nThe following describes each of the modes more specifically. Mode 0 carries out a spatial prediction toward a vertical direction with the use of the pixels A through D. Mode 1 carries out the spatial prediction toward a horizontal direction with the use of the pixels I through L. Prediction Mode 2 carries out a DC prediction with the use of the pixels A through D and the pixels I through L. Mode 3 carries out the spatial prediction toward a diagonal down left direction with the use of the pixels A through H. Mode 4 carries out the spatial prediction toward a diagonal down right direction with the use of the pixels A through D and the pixels I through M. Mode 5 carries out the spatial prediction toward a vertical right direction with the use of the pixels A through D, the pixels I through K, and the pixel M. Mode 6 carries out the spatial prediction toward a horizontal down direction with the use of the pixels A through C and the pixels I through M. Mode 7 carries out the spatial direction toward a vertical left direction with the use of the pixels A through E. Mode 8 carries out the spatial prediction toward a horizontal up direction with the use of the pixels I, J, K, and L. The intra prediction section 310 generates prediction pixels by use of a prediction method corresponding to any one of the modes.\n(Prediction Mode Determining Section 311)\nThe prediction mode determining section 311 determines which one of the plurality of prediction modes shown in FIG. 66 is to be used for prediction of the target block to be encoded, based on an inputted original image of the target block to be encoded and a locally decoded image received from the memory 9, and supplies information on a determined prediction mode (hereinafter referred to as prediction mode information) to the intra prediction section 310 and the entropy encoding section 304.\nGenerally, a method for evaluating a prediction residual cost (hereinafter referred to as a prediction residual cost minimization method) or a rate distortion optimization method is used for determining a prediction mode to be used.\nThe prediction residual cost minimization method is a method in which (i) a degree of similarity (hereinafter referred to as a prediction residual cost) is found between the inputted original image of the target block to be encoded and a predicted image corresponding to each of the prediction modes, generated on the basis of a locally decoded image received from the memory 9, and then (ii) a prediction mode having the smallest prediction residual cost is selected among the prediction modes. A measure S of the prediction residual cost is represented by any one of the followings: a sum of absolute values of prediction residual data; a square sum of prediction residual data; a sum of absolute values of transform coefficients of prediction residual data; and a square sum of transform coefficients of prediction residual data. These values are calculated by use of the following equations (1) through (4).\n                    Math        .                                  ⁢        1                                                            S        =                              ∑                          i              ,                              j                ∈                block                                              ⁢                                                                f                ⁡                                  (                                                            x                      +                      i                                        ,                                          y                      +                      j                                                        )                                            -                              p                ⁡                                  (                                                            x                      +                      i                                        ,                                          y                      +                      j                                                        )                                                                                                    (        1        )                                S        =                              ∑                          i              ,                              j                ∈                block                                              ⁢                                    {                                                f                  ⁡                                      (                                                                  x                        +                        i                                            ,                                              y                        +                        j                                                              )                                                  -                                  p                  ⁡                                      (                                                                  x                        +                        i                                            ,                                              y                        +                        j                                                              )                                                              }                        2                                              (        2        )                                S        =                              ∑                          i              ,                              j                ∈                block                                              ⁢                                                T              ⁢                              {                                                      f                    ⁡                                          (                                                                        x                          +                          i                                                ,                                                  y                          +                          j                                                                    )                                                        -                                      p                    ⁡                                          (                                                                        x                          +                          i                                                ,                                                  y                          +                          j                                                                    )                                                                      }                                                                                    (        3        )                                S        =                              ∑                          i              ,                              j                ∈                block                                              ⁢                      T            ⁢                                          {                                                      f                    ⁡                                          (                                                                        x                          +                          i                                                ,                                                  y                          +                          j                                                                    )                                                        -                                      p                    ⁡                                          (                                                                        x                          +                          i                                                ,                                                  y                          +                          j                                                                    )                                                                      }                            2                                                          (        4        )            \nIn the equations (1) through (4), f(x, y) indicates an original image, p(x, y) indicates a predicted image, x and y indicate a target block to be encoded, and i and j indicate a position of a pixel in the target block to be encoded. T{ } indicates an orthogonal transformation operation, such as discrete cosine transform, discrete sine transform, and Hadamard transform.\nIn the rate distortion optimization method, a prediction mode is selected according to the following steps. Initially, predicted images corresponding to the respective prediction modes are generated on the basis of a locally decoded image received from the memory 9. Then, prediction residual data is found from each of the predicted images thus generated and the inputted original image of the target block to be encoded. Subsequently, the prediction residual data thus obtained is temporarily encoded. Then, (a) a prediction error D between the original image of the target block and a decoded image of the target block and (b) an encoding amount R necessary for encoding the target block are calculated with respect to each of the prediction modes. Finally, a prediction mode that minimizes an encoding cost J calculated according to the following equation (5) based on D and R, is selected from the prediction modes.\nMath. 2J(mode|q)=D(mode|q)+λ(q)·R(mode|q)  (5)\nThe prediction error D is an error between prediction residual data that has not been quantized and prediction residual data that has quantized. The encoding amount R is a total of an encoding amount of prediction residual data and an encoding amount of prediction mode information.\nIn the equation (5), “mode” indicates a prediction mode, and q indicates a quantization parameter. Further, λ is a weighting factor dependent on the quantization parameter q, and is generally calculated according to the following equation (6):\nMath. 3λ(q)=0.85×2(q−12)/3  (6)\n(Encoding Process of Image Encoding Device 300)\nThe following describes how the image encoding device 300 operates.\nInitially, a target image to be encoded (hereinafter referred to as a target block to be encoded), which is divided into blocks by a predetermined block size (M×M pixel block), is supplied to the image encoding device 300, as an input image.\nThe prediction mode determining section 311 determines a prediction mode to be used for prediction for the target block to be encoded, in accordance with the prediction residual cost minimization method or the rate distortion optimization method, based on the target block to be encoded and a locally decoded image of an adjacent block that has been already encoded, which locally decoded image is received from the memory 9. The prediction mode determining section 311 then supplies prediction mode information to the intra prediction section 310 and the entropy encoding section 304.\nThe intra prediction section 310 carries out intra prediction indicated by the prediction mode information thus received, based on the locally decoded image of the adjacent block that has been encoded, which locally decoded image is received from the memory 9, so as to generate a predicted image (M×M pixel block) of the target block to be encoded. The intra prediction section 310 then supplies the predicted image to the difference computing section 1.\nSubsequently, the difference computing section 1 calculates prediction residual data that is a difference between the target block to be encoded and the predicted image thus generated, and supplies the prediction residual data to the orthogonal transformation section 2.\nThe prediction residual data from the difference computing section 1 is supplied to the orthogonal transformation section 2 and then to the quantization section 3 so as to be subjected to orthogonal transformation and then quantization. The prediction residual data is then supplied to the entropy encoding section 304 and to the inverse quantization section 6.\nThe prediction residual data thus subjected to the orthogonal transformation and the quantization is supplied to the inverse quantization section 6 and then to the inverse orthogonal transformation section 7 so as to be subjected to inverse quantization and inverse orthogonal transformation, respectively. The prediction residual data is then supplied to the addition computing section 8.\nThe addition computing section 8 combines the prediction residual data thus subjected to the inverse quantization and the inverse orthogonal transformation, with the predicted image corresponding to the prediction mode applied to the target block, so as to synthesize a locally decoded image (M×M pixel block) of the target block. The locally decoded image thus synthesized is then supplied to the memory 9.\nThe locally decoded image of the target block, which is supplied from the addition computing section 8, is stored in the memory 9, and is used for intra prediction of an adjacent block to be encoded subsequently to the target block.\nThe entropy encoding section 304 carries out an encoding process, such as variable length coding, with respect to encode parameters, such as the prediction mode information, received from the prediction mode determining section 311 and the prediction residual data subjected to the orthogonal transformation and the quantization. Subsequently, the entropy encoding section 304 outputs encoded data of the target block.\nThe image encoding device 300 repeats the above process with respect to all target blocks to be encoded, which constitute a target image to be encoded.\nAs described above, the image encoding device 300 selects a corresponding one of the 9 types of predetermined prediction methods from Modes 0 through 8 illustrated in FIG. 66, depending on a characteristic of an image. This can attain high encoding efficiency. More specifically, the DC prediction of Mode 2 is effective for prediction for a planar image region. Further, Modes 0, 1 , and 3 through 8 are effective for prediction for an image including an edge in a specific direction.\n(a) and (b) of FIG. 67 illustrate concrete examples of the intra prediction in the image encoding device 300. (a) of FIG. 67 illustrates an example in which prediction residual data F43 is obtained as a difference between (i) an image F41 having vertical stripes in which edges are formed in a vertical direction and (ii) a predicted image F42 generated, based on a decoded image, according to a vertical direction prediction (Mode 0 in FIG. 66). Further, (b) of FIG. 67 illustrates an example in which prediction residual data F46 is obtained as a difference between (i) an image F44 having diagonal stripes in which edges are formed in a diagonal right direction and (ii) a predicted image F45 generated, based on a decoded image, according to a diagonal down right direction prediction (Mode 4 in FIG. 66).\nIn the cases of (a) and (b) of FIG. 67, the predicted images F42 and F45 respectively reproduce patterned parts in the images F41 and F44, which are the original images. Therefore, when a difference between the original image and the predicted image is calculated, the patterned part in the original image and the patterned part in the predicted image cancel each other out. As a result, pixels, in the original image, which cannot be subtracted by the predicted image remains as prediction residual data. In (a) and (b) of FIG. 67, a prediction direction is identical with a direction of the pattern in the original image. Consequently, it is possible to effectively reduce the prediction residual data.\n(Predictive Encoding Method of Prediction Mode Information)\nThe following deals with a predictive encoding method of prediction mode information.\nThe intra prediction of a 4×4 pixel block or a 8×8 pixel block in the image encoding device using a spatial correlation, disclosed in Non Patent Literature 1, uses the 9 types of prediction modes as illustrated in FIG. 68. In Non Patent Literature 1, when prediction mode information of a target block (B1) is encoded, an estimation value is determined in such a manner that (a) a value of a prediction mode applied to a block B2, which is adjacently positioned on the left side of the block B1, and (b) a value of a prediction mode applied to a block B4, which is adjacently positioned on the upper side of the block B1, are compared, and the prediction mode having a smaller value between them is taken as the estimation value. Then, the estimation value is compared with a value of a prediction mode of the block B1. If the values are identical, a flag “1” is encoded. In the other cases, a flag “0” and relative prediction mode information indicative of which one of the remaining 8 types of prediction modes (i.e., the prediction modes except for the estimation value) corresponds to the prediction mode of the block B1 are encoded.\nFurther, in Patent Literature 1, when the prediction mode information of the target block is encoded, a prediction mode applied to the target block is estimated by use of an estimation value 1 determined by the after-mentioned first prediction means and an estimation value 2 determined by the after-mentioned second prediction means.\nThe first prediction means determines, as the estimation value 1, a prediction mode that is most used among prediction modes applied to (i) blocks in an encoded region C2 (for example, a macro block, a slice) which is positioned on the left side of a target block C1, or alternatively (ii) blocks in an encoded region C3 which is positioned on the upper side of the target block C1. Further, the second prediction means uses a decoded image in an L-shaped region B5 (a shaded part with diagonal down-right lines) on encoded blocks B2, B3 and B4 which are adjacent to a target block B1, as illustrated in FIG. 69. More specifically, predicted images corresponding to 9 types of intra prediction methods shown in FIG. 68 are generated with respect to a decoded image of an L-shaped region B6 (a shaded part with diagonal down-left lines), with the use of the decoded image of the L-shaped region B5. Then, a prediction mode whose difference between the decoded image of the L-shaped region B6 and the predicted image is smallest among those prediction modes is regarded as an estimation value 2.\n(Arrangement of Image Decoding Device 350)\nFinally, a conventional image decoding device is described as below with reference to FIG. 71. FIG. 71 is a block diagram showing an arrangement of a conventional image decoding device 350. The image decoding device 350 includes an inverse quantization section 6, an inverse orthogonal transformation section 7, an addition computing section 8, a memory 9, an entropy decoding section 305, and an intra prediction section 310 (FIG. 71). Among these constituent components, the inverse quantization section 6, the inverse orthogonal transformation section 7, the addition computing section 8, the memory 9, and the intra prediction section 310 have been described above. Therefore, the following deals with only the entropy decoding section 305.\n(Entropy Decoding Section 305)\nThe entropy decoding section 305 carries out a decoding process (for example, a variable length decoding process) with respect to encoded data indicative of prediction residual data of a target block to be decoded and encode parameters, such as prediction mode information, of the target block to be decoded.\n(Decoding Process of Image Decoding Device 350)\nThe following describes a decoding process of an image in the image decoding device 350. Initially, the entropy decoding section 305 carries out entropy decoding with respect to received encoded data of the target block to be decoded, and outputs prediction residual data of the target block to be decoded and encode parameters, such as prediction mode information, of the target block to be decoded. The prediction residual data thus decoded is supplied to the inverse quantization section 6 and then to the inverse orthogonal transformation section 7 so as to be subjected to inverse quantization and inverse orthogonal transformation, respectively.\nThen, the intra prediction section 310 generates a predicted image (M×M pixel block) of the target block to be decoded with the use of a locally decoded image, stored in the memory 9, of an adjacent block that has been already decoded.\nSubsequently, the addition computing section 8 adds the prediction residual data thus subjected to the inverse quantization and the inverse orthogonal transformation, to the predicted image generated by the intra prediction section 310, so as to generate a decoded image of the target block. The memory 9 stores the decoded image (M×M pixel block) thus generated therein. The stored decoded image is used for intra prediction of an adjacent block(s). The image decoding device 350 repeats the above process with respect to all target blocks to be decoded, which constitute a target image to be decoded."}
{"text": "This invention relates generally to devices which support and hold equipment and more particularly to a foot for such devices suitable for holding geomatic equipment and accessories.\nAlthough the art of surveying is old, recent improvements in equipment, such as automatic electronic total surveying stations, have increased the accuracy of the surveying instruments far beyond what was possible previously. Where accepted accuracy for equipment in the not too distant past was 20 seconds of a degree, accuracy is now commonly required to be 1 second of a degree. At these levels of precision, the instruments and many of their accessories must be supported in a manner which is consistent and stable. Accordingly, there is a need for closer examination of support structure used for surveying equipment so that improvements in instrumentation will not be lessened by inadequate supports. Further, the cost of surveying equipment makes it highly undesirable to damage it by failure of the supports because of instability.\nPole supports used in the surveying field typically take the form of bipods or tripods, but conceivably any number of legs could be provided. Conventionally, it has been necessary to provide completely separate inventories of bipods and tripods to meet the needs of different customers and for different applications. Legs of these supports are usually freely pivotable about a head of the support to swing toward each other for ease of carrying the support when not in use. A support having three or more legs can support equipment or other items above and out of contact with the ground. Typically, bipods are used to support surveying equipment such as a surveying pole, which also contacts the ground, but tripods may also be used to support surveying poles. The stability of such supports can be compromised in situations where there are forces (e.g., wind and loads from the surveying equipment) acting on the support and tending to tip the support over. In that case, one of the legs is subjected to an upward force component tending to raise the leg up as the support pivots about a point(s) of contact of the other leg(s) with the ground.\nPresently, there is little or no resistance provided by the legs to such upward forces other than the weight of the leg. Even if there are stakes associated with the legs that penetrate the ground, the upward force tends to be directed right along the long axis of the stakes so that they are lifted out of the ground. The problem is compounded by the fact that the legs conventionally are mounted for pivoting freely about a horizontal axis. Even if the support does not tip over, if one leg loses contact with the ground or the frictional resistance to pivoting movement of the leg becomes sufficiently small, the leg will swing in toward the center of the support. The support is unlikely to regain stability when this happens. Should the support move back toward its original position, the leg having pivoted inwardly is no longer in position to engage the ground in a stable position. The center of gravity of the combined support and surveying equipment may lie outward of the point where the leg re-contacts the ground so that the entire unit topples over in the opposite direction from the initial tipping movement. Alternatively as the leg pivots inwardly, the weight of the leg exerts a smaller torque resisting the tipping motion so that the motion may continue causing the unit to collapse. If instability of the support causes it to move, even where it does not fall over, accuracy is compromised.\nThe legs of the support are most often telescoping so that their length may be adjusted as needed so that, for instance, the support can be set up on uneven terrain or the ends of the legs contact the ground at different leg angles. In the situation where the support is used with a surveying pole, the legs are manipulated so that the pole is precisely located over a selected point. When the telescoping legs are extended, it is necessary that they be locked at a specific length. Locking systems work most commonly by greatly increasing the frictional engagement of one leg section with the other. These systems suffer from certain drawbacks. Some locking systems may have a component of motion which is parallel to the lengthwise extent of the leg when engaged. This movement can cause the length of the leg to change during engagement of the locking system and result in loss of precision in the position of a surveying pole held by the support. The leg sections may be difficult to move relative to each other because of residual frictional engagement of the leg sections caused by the locking system, even when the locking system is released. In that situation movement may be jerky, making it difficult to extend or retract the leg to a precise length. Moreover, it is conventionally necessary to manually hold the locking device to keep the leg sections unlocked. This makes manipulation of the legs difficult, particularly where there are more than two legs. Conventionally, the leg sections are assembled at the factory and cannot be readily taken apart. Accordingly, cleaning and repair in the field is difficult.\nIt may be difficult to position the surveying pole in the correct location when a support is attached to the pole and the legs of the support are in a stowed position. The legs can obstruct the surveyor's view of the bottom point of the pole. It is desirable for the legs to be capable of retracting to a short length to reduce the obstruction. However, any items attached to the lower ends of the legs to facilitate securing the ends in the ground project down below the lower ends, increasing the effective retracted length of the leg and making it more difficult to remove the items as an obstruction by retracting the length of the legs.\nFeatures of convenience are usually neglected in a conventional support. At most, a strap is attached to one of the legs which can be wrapped around the legs and secured to hold the legs together in their collapsed configuration. Manipulation of the strap to hold or release the legs can be cumbersome and inconvenient. Little has been done to provide an easy, way of securing the legs in the collapsed position. The other surveying tools and equipment which will be needed are carried separately. It is known to provide dedicated supports to mount items on a leg of a support or on a surveying pole. These suffer from the problem of causing point or line loads on the leg or pole which can cause damage and/or reduce the accuracy of the support or surveying pole. The problem is particularly acute where lighter weight materials are used to construct a surveying pole. Still further, there is no connecting system which would readily allow attachment of different articles to the same mount associated with the leg or surveying pole. Moreover, there is presently no single connection system for connecting articles to two separate objects. Articles which are not presently needed in surveying must be held, or loosely stored in a case or bag.\nAlthough the present invention has particular application for use in supporting surveying and geographic positioning equipment, it is not limited to such applications. As used herein “geomatic” is intended to encompass surveying and geographic positioning. The invention is envisioned as being useful to support equipment and other items having no relation to surveying or geographic positioning."}
{"text": "Currently, in situations where Global Position System (GPS location information is not available, a user typically relies on dead reckoning localization (e.g., using inertial measurements from an inertial measurement unit (IMU)). Such localization analysis, however, is subject to drift and error accumulation. Yaw, for example, is typically calculated using a compass, which can be unreliable due to variations of the earth's magnetic field direction on the surface of the earth, meaning yaw measurements using a compass can be inaccurate by many degrees. Other methods for calculating yaw include measuring celestial features such as the Sun, moon and star positions. These methods can be accurate (less than 1 degree of error), but are subject to reduced availability due to cloud cover, the Sun being out of the field of view, stars not being visible during the daytime, and the like.\nAccording to theory, the observed polarization at any position in the sky depends on the Sun and the sensor platform positions, as well as the sensor pointing direction, where “sensor pointing direction” is the center point of the field of view of the sensor, also known as the target point. The target point, sensor platform position, and sun position together define a plane. Given the Sun's position, which is a function of the time of day, and polarization measurements at one or more unique pointing directions, the sensor absolute position and orientation may be derived. As used herein, “orientation” generally refers to roll, pitch and yaw. “Position” generally refers to latitude and longitude.\nA method according to the present disclosure calculates orientation and position parameters using a sky polarimeter that takes polarized images of multiple simultaneous target points in the sky. The orientation and position parameters can be useful to a navigating vehicle (especially if GPS is denied, spoofed, or unavailable), and can work in all types of vehicles (including ground, air and naval vehicles). The orientation and position parameters can also be useful to target locating systems such as far target locators and surveying equipment. The method can provide 0.1 degree yaw accuracy. Further, while the method is typically applied during daylight hours, it is conceivable that the method could be executed at night with some accuracy using the moon instead of the sun.\nA system according to an exemplary embodiment of the present disclosure comprises an imaging sensor, polarization state analyzer, optics, mechanical housing, memory and logic circuitry, IMU, GPS, clock, and embedded software that determine the orientation and position parameters. A method according to an exemplary embodiment of the present disclosure comprising using polarization images and prior position/orientation/time data from the GPS, IMU and clock, respectively, to determine expected Sun azimuth and elevation, comparing this expected Sun position to the measured sky polarization pattern, and then filtering to calculate a better orientation and position estimate of the desired object. This localization estimate can be provided in any number of interfaces to a navigation system, a user, a display, or a target locator."}
{"text": "Improvements in magnetic materials and commutation techniques have enabled dynamoelectric machines to be significantly reduced in size while developing as much or more power than their larger, bulky predecessors. In some instances, the size reductions have resulted in machine structures having less than desirable rigidity. For example, some high torque, low speed motors are constructed with many poles at a large radius and have limited stiffness in the radial direction. Some of these motors may have an air gap diameter of about eight inches and a lamination stack height of about one and one-half inches. In order to maintain uniform circularity of the lamination stack so as to maintain a uniform air gap between stator and rotor, it is necessary to support the stator and rotor in a rigid frame. However, it is desirable to minimize the mass of frame structure necessary to provide such support.\nIn switched reluctance (SR) motors or in single phase motors, rigidity is particularly important since such motors generate torque pulsations which produce noise unless the magnetic structure is mechanically isolated from the base structure by a torsionally soft mount. Accordingly, it is desirable to provide a mounting structure for such dynamoelectric machines of relatively low mass, suitable for low cost mass production and which provides isolation of radial and torsional vibrations from the base structure."}
{"text": "This invention relates to synthetic organic chemistry. Specifically, the invention relates to a process for preparing intermediates useful in the syntheses of valuable antifolate compounds.\nCompounds known to have antifolate activity are well recognized as chemotherapeutic agents for the treatment of cancer. A series of N-(6-amino-(pyrrolo(2,3-d)pyrimidin-3-ylacyl)-glutamic acid derivatives of formula V: \nwhere\nY is CHxe2x95x90CH, O, or S;\nR3 is hydrogen, chloro, or fluoro;\nR4 is hydroxy, a carboxy protecting group, or NHCH*(C(D)R5)CH2CH2C(O)R5;\nR5 is hydrogen or a carboxy protecting group;\nR6 is hydrogen or an amino protecting group;\nR7 is hydroxy or amino; and the configuration about the carbon atom designated * is S; and the pharmaceutical salts thereof were disclosed as antifolates or intermediates to antifolates in U.S. Pat. Nos. 4,684,653 and 4,882,334.\nA key step in the synthesis of the compounds of formula V, disclosed in U.S. \"\"334 and \"\"653, is the hydrogenation of compounds of formula VI: \nwhere Z1 and Z2 are both hydrogen or taken together form a bond; R4xe2x80x2 is a carboxy protecting group or NHC*H(C(O)R5xe2x80x2)CH2CH2C(O)R5xe2x80x2;\nR5xe2x80x2 is a carboxy protecting group; and\nR6xe2x80x2 is an amino protecting group;\nproviding the isomeric mixture of compounds of formula V(a): \nThe resulting compound of formula V(a) can optionally have its protecting groups removed to give an isomeric mixture of the compounds of formula V. U.S. \"\"334 and \"\"653 further taught that the individual diastereomers of formula V could be separated mechanically by chromatography or preferably the individual diastereomers could be separated by forming diastereomeric salts with chiral acids, such as camphorsulfonic acid, followed by selective crystallization of one of the diastereomers.\nU.S. \"\"334 and \"\"653 taught that compounds of formula VI can be obtained by first coupling a compound of formula VII with a compound of formula VIII: \nwhere X is bromo or iodo; in the presence of a palladium/trisubstituted phosphine catalyst of the type described by Sakamoto in Synthesis, 1983, 312 et. seq.\nThe synthesis outlined above suffers in many respects. On an industrial scale, use of a noble metal catalyst is expensive, leads to purification and environmental issues, and can be erratic due to varying amounts of the precious metal that is in the correct oxidation state/complex form for catalysis. Furthermore, if the preferred crystallization separation procedure taught above is followed, after isolating a diastereomer by filtration, the filtrate will contain mixtures of the two diastereomers. This filtrate is often not amenable to further separation by crystallization, and thus separation efficiency suffers without resorting to an undesired chromatographic separation. In certain cases, e.g., where Y is S, R3 is hydrogen, R4 is NHC*H(C(O)R5)CH2CH2C(O)R5, and R7 is hydroxy, as much as 80% of the desired isomer (the one with greater antifolate activity) could be found in the filtrate/fractions.\nAn improvement over the prior art would not rely on precious metal catalysis to produce the compounds of formula VI and would increase the absolute yields of the desired diastereomer of formula VI from mixtures containing both diastereomers by crystallization.\nThe present invention relates to a compound of formula III: \nwhere:\nY is CHxe2x95x90CH, O, or S;\nR is C1-C6 alkyl;\nR1 and R2 are independently C1-C6 alkyl or taken together with the nitrogen to which they are attached form a heterocycle;\nR3 is hydrogen, chloro, or fluoro;\nR4 is hydroxy, a carboxy protecting group, or NHC*H(C(O)R5)CH2CH2C(O)R5 where the configuration about the carbon atom designated * is S; and\nR5 is hydrogen or a carboxy protecting group; or a salt or solvate thereof.\nThe present invention also relates to a compound of formula IV: \nwhere:\nR6 is hydrogen or an amino protecting group; and\nR7 is hydroxy or amino; or a salt or solvate thereof.\nMoreover, the present invention relates to a process for preparing compounds of formula IV: \nwhich includes reacting a compound of formula III(a): \nwhere:\nR4xe2x80x2 is a carboxy protecting group or NHC*H(C(O)R5xe2x80x2)CH2CH2C(O)R5xe2x80x2 where the configuration about the carbon atom designated * is S; and\nR5xe2x80x2 is a carboxy protecting group;\nwith 2,4-diamino-6-hydroxypyrimidine in the presence of a suitable acid and solvent.\nFurthermore, the present invention also relates to a process for preparing a compound of formula IV, or a salt or solvate thereof, which includes reacting a compound of formula V(b): \nor a salt or solvate thereof, with an oxidizing reagent in the presence of a suitable solvent.\nIn the general formulae of the present document, the general chemical terms have their usual meanings. For example, the term xe2x80x9cC1-C4 alkylxe2x80x9d refers to methyl, ethyl, propyl, isopropyl, cyclopropyl, n-butyl, s-butyl, t-butyl, and cyclobutyl. The term xe2x80x9cC1-C6 alkylxe2x80x9d encompasses those listed for C1-C4 alkyl in addition to aliphatic, straight, branched, or cyclic, monovalent moieties having five or six carbon atoms and includes, but is not limited to, pentyl, cyclopentyl, hexyl, cyclohexyl, 2-methylpentyl, and the like. The term xe2x80x9cC1-C4 alkoxyxe2x80x9d refers to a C1-C4 alkyl group attached through an oxygen atom.\nThe term xe2x80x9chaloxe2x80x9d or xe2x80x9chalidexe2x80x9d refers to chloro, bromo, or iodo.\nThe term xe2x80x9cheterocyclexe2x80x9d refers to a 5 or 6 membered saturated, partially unsaturated, or aromatic heterocyclic ring which contains a nitrogen atom and may optionally contain an additional heteroatom selected from N, S, or O.\nThe term xe2x80x9ccarboxy protecting groupxe2x80x9d refers to a substituent of a carbonyl that is commonly employed to block or protect the carboxy functionality while reactions are carried out on other functional groups on the compound.\nThis substituent, when taken with the carbonyl to which it is attached, may form an ester, e.g., C1-C6 alkyl, substituted C1-C6 alkyl, C2-C6 alkenyl, substituted C2-C6 alkenyl, benzyl, substituted benzyl, benzhydryl, substituted benzhydryl, trityl, substituted trityl, and trialkylsilyl ester. The exact species of carboxy protecting group is not critical so long as the derivatized carboxy group is stable to the conditions of subsequent reaction(s) and can be removed at the appropriate point without disrupting the remainder of the molecule. When R4 contains a carboxy protecting group, the protecting group is preferably C1-C4 alkoxy or benzyloxy. The most preferred protecting groups are methoxy, ethoxy, and benzyloxy. A carboxy protecting group xe2x80x9cremovable by catalytic hydrogenationxe2x80x9d includes, for example, benzyl protecting groups. Other examples of these groups are described in T. W. Greene, xe2x80x9cProtective Groups in Organic Synthesis,xe2x80x9d John Wiley and Sons, New York, N.Y., (2nd ed., 1991), (hereafter referred to as Greene) chapter 5.\nThe term xe2x80x9cC2-C6 alkenylxe2x80x9d refers to a mono-unsaturated, monovalent, hydrocarbon moiety containing from 2 to 6 carbon atoms which may be in a branched or straight chain configuration. The term is exemplified by moieties such as, but not limited to, ethylenyl, propylenyl, allyl, butylenyl, and pentylenyl.\nThe terms xe2x80x9csubstituted C1-C6 alkylxe2x80x9d and xe2x80x9csubstituted C2-C6 alkenylxe2x80x9d refer to a C1-C6 alkyl and C2-C6 alkenyl group respectively substituted from 1 to 3 times independently with a halo, phenyl, tri(C1-C4 alkyl)silyl, or a substituted phenylsulfonyl group.\nThe terms xe2x80x9csubstituted benzylxe2x80x9d, xe2x80x9csubstituted benzhydrylxe2x80x9d, and xe2x80x9csubstituted tritylxe2x80x9d refers to a benzyl, benzhydryl, and trityl group, respectively, substituted from 1 to 5 times independently with a nitro, C1-C4 alkoxy, C1-C6 alkyl, or a hydroxy(C1-C6 alkyl) group. These substitutions will only occur in a sterically feasible manner such that the moiety is chemically stable.\nThe term xe2x80x9ctrialkylsilylxe2x80x9d refers to a monovalent silyl group substituted 3 times independently with a C1-C6 alkyl group.\nThe term xe2x80x9csubstituted phenylsulfonylxe2x80x9d refers to a henylsulfonyl group where the phenyl moiety is para ubstituted with a C1-C6 alkyl, nitro, or a halo group.\nThe term xe2x80x9camino protecting groupxe2x80x9d as used in the specification refers to a substituent of the amino group commonly employed to block or protect the amino functionality while reacting other functional groups on the compound. The amino protecting group, when taken with the nitrogen to which it is attached, can form a cyclic imide, e.g., phthalimido and tetrachlorophthalimido; a carbamate, e.g., methyl, ethyl, and 9-fluoroenylmethylcarbamate; or an amide, e.g., N-formyl and N-acetylamide. The exact genus and species of amino protecting group employed is not critical so long as the derivatized amino group is stable to the condition of subsequent reaction(s) on other positions of the intermediate molecule and can be selectively removed at the appropriate point without disrupting the remainder of the molecule including any other amino protecting group(s). In general, amino protecting groups removable by acid hydrolysis, i.e., those that are acid labile, are preferred. Thus, a preferred amino protecting groups is 2,2-dimethyl-1-oxopropyl. Further examples of groups and methods referred to by the above terms are described in Greene at chapter 7.\nThe term xe2x80x9cpharmaceutical saltxe2x80x9d and xe2x80x9csaltxe2x80x9d as used herein, refers to salts prepared by reaction of the compounds of the present invention with a mineral or organic acid (e.g., hydrochloric, hydrobromic, hydroiodic, or p-toluenesulfonic acid) or an inorganic base (e.g., sodium, potassium, lithium, and magnesium hydroxide, carbonate, or bicarbonate). Such salts are known as acid addition and base addition salts. For further exemplification of these salts, see, e.g., Berge, S. M, Bighley, L. D., and Monkhouse, D. C., J. Pharm. Sci., 66, 1, 1977.\nThe term xe2x80x9csolvatexe2x80x9d represents an aggregate that comprises one or more molecules of a solute, such as a formula III or IV compound, with one or more molecules of solvent.\nThe term xe2x80x9csuitable solventxe2x80x9d refers to any solvent, or mixture of solvents, inert to the ongoing reaction that sufficiently solubilizes the reactants to afford a medium within which to effect the desired reaction.\nThe term xe2x80x9csuitable acidxe2x80x9d refers to an acid whose Ka is low enough to effect the desired reaction without significantly effecting any undesired reactions.\nThe term xe2x80x9coxidizing reagentxe2x80x9d refers to a reagent whose oxidation potential is high enough to effect the desired reaction without significantly effecting any undesired reactions. Suitable oxidants include metals such as nickel, palladium, platinum, and the like; metals on solid supports such as palladium or platinum on carbon, and the like; metal complexes such as mercury(II), manganese dioxide, or coppor(II) acetate, and benzoquinone based oxidants such as 2,3-dichloro-5,6-dicyano-1,4-benzoquinone (DDQ) and tetrachloro-1,4-benzoquinone (chloranil), and the like.\nThe term xe2x80x9cthermodynamic basexe2x80x9d refers to a base which provides a reversible deprotonation of an acidic substrate or is a proton trap for those protons that may be produced as byproducts of a given reaction, and is reactive enough to effect the desired reaction without significantly effecting any undesired reactions. Examples of thermodynamic bases include, but are not limited to, acetates, acetate dihydrates, carbonates, bicarbonates, C1-C4 alkoxides, and hydroxides (e.g., lithium, sodium, or potassium acetate, acetate dihydrate, carbonate, bicarbonate, C1-C4 alkxoxide, or hydroxide), tri(C1-C4 alkyl)amines, or aromatic nitrogen containing heterocycles (e.g., imidazole and pyridine).\nThe compounds of formula III may be prepared from compounds of formula I and II as illustrated in Scheme 1 below where Lg is chloro, bromo, iodo, OSO2Me, OSO2-phenyl, or OSO2-p-toluenyl and R, R1, R2, R3, R4, R4xe2x80x2, and Y are as defined above. \nCompounds of formula I, may be added to compounds of formula II dissolved or suspended in a suitable solvent, in the presence of a thermodyanamic base, to form the compounds of formula III(a). A preferred and convenient solvent is dichloromethane. A preferred and convenient base is triethylamine. A single equivalent of base and compound of formula I, relative to the compound of formula II, is preferably employed but slight excesses on the order of 0.01 to 0.1 equivalents are tolerable. The reaction may be performed between xe2x88x9278xc2x0 C. and ambient temperature but is preferably performed between xe2x88x9225xc2x0 C. and xe2x88x9220xc2x0 C. The reaction is typically complete in from 30 minutes to 18 hours but when performed at the preferred temperature, it is complete in from 1 to 3 hours. A preferred halide in compounds of formula I is chloride. R is preferably C1-C4 alkyl, especially methyl or ethyl. R1 and R2 are preferably C1-C4 alkyl but it is especially preferred when both are either methyl or ethyl. It is preferred that R4xe2x80x2 is a carboxy protecting group where that protecting group is C1-C4 alkoxy, especially methoxy or ethoxy, or one capable of being removed by catalytic hydrogenation, e.g., benzyloxy (as in Scheme 3 below). Throughout this specification, R3 is preferably hydrogen and Y is preferably CHxe2x95x90CH or S.\nAlthough the resulting compound of formula III(a) may have its carboxy protecting group removed as taught in Greene, for the purposes of conducting the overall process of Schemes 1-3, the carboxy protecting is preferably left intact when proceeding to the reaction(s) of Scheme 2.\nCompounds of formula IV may be prepared from compounds of formula III(a) by the novel process illustrated in Scheme 2 below where R, R1, R2, R3, R4, R4xe2x80x2, R6, R7, and Y are as defined above. \nA tautomeric mixture of 2,4-diamino-6-hydroxypyrimidine or 2,4-diaminopyrimidin-6-one (hereafter referred to as 2,4-diamino-6-hydroxypyrimidine) may be added to a compound of formula III(a), dissolved or suspended in a suitable solvent, in the presence of a suitable acid, to provide the compounds of formula IV(a). A convenenient and preferred solvent is an approximately 1:1 (v:v) mixture of acetonitrile and water. The ratio of acetonitrile to water is not critical but it is preferred that the ratio is amenable to forming a solution when the reactants are all initially combined. A convenient and preferred acid is acetic acid. The acid is typically employed in molar excess. For example, about 2 to about 4 equivalents, relative to the compound of formula III(a), is generally employed while 3 equivalents are typically preferred. The number of equivalents of 2,4-diamino-6-hydroxypyrimidine employed relative to the compound of formula III(a) is not critical but about 1 to about 2 equivalents are preferred. An even more preferred amount is about 1 to about 1.5 with about 1 to about 1.1 equivalents most preferred. The reaction may be performed at temperatures ranging from room temperature to the reflux temperature of the mixture but is preferably performed at the reflux temperature of the mixture. Furthermore, the reaction may take from 12 to about 48 hours depending on the temperature of the reaction. When the reaction is performed at the reflux temperature of the mixture, it is typically substantially complete in about 18 hours. R is preferably C1-C4 alkyl, especially methyl or ethyl. R1 and R2 are preferably C1-C4 alkyl but it is especially preferred when both are either methyl or ethyl. R3 is preferably hydrogen. As stated previously, it is preferred that R4xe2x80x2 is a carboxy protecting group where that protecting group is a C1-C4 alkoxy group, especially methoxy or ethoxy, or one capable of being removed by catalytic hydrogenation (as in Scheme 3).\nThe compounds of formula IV where R7 is an amino group may be prepared from the compounds of formula IV(a) as taught in the previously incorporated by reference U.S. Pat. No. 4,882,334 but it is preferred that R7 is hydroxy.\nThe compounds of formula IV where R6 is an amino protecting group may be prepared from compounds of formula IV(a) as taught in Greene or as discussed in Preparation 8 below. Furthermore, it is necessary that an amino protecting group, preferably one removable by acid hydrolysis as in Scheme 5 below, e.g., 2,2-dimethyl-1-oxopropyl be present at R6 or that the amino group be protonated before proceeding to the hydrogenation described in Scheme 3 below. Moreover, although the compounds of formula IV(a) may have their carboxy protecting groups removed as taught in Greene, for the purposes of conducting the overall process of Schemes 2-3, it is preferred that the carboxy protecting group is left intact when proceeding to the reaction of Scheme 3.\nAn isomeric mixture of compounds of formula V(b) may be prepared from compounds of formula IV(b) as illustrated in Scheme 3 below where R3, R4, R6xe2x80x2, R7, and Y are as defined above. \nCompounds of formula IV(b), prepared as described in Schemes 1 and 2, may be hydrogenated substantially as described in U.S. Pat. Nos. 4,684,653 and 4,882,334, the teachings of each are hereby incorporated by reference. For facilitation of cross reference, the compounds of formula IV(b) in the present invention correspond to the compounds of formula III in U.S. \"\"653 and the compounds of formula II in U.S. \"\"334. It is preferred that the amino protecting group at R6xe2x80x2 is not removed as taught in Greene but left intact when continuing to the separation procedures discussed in Schemes 4 and 5 below.\nIf the process of Scheme 3 is performed with the compounds of formula IV(b) with the preferred group at R4, i.e., a carboxy protecting group, and that carboxy protecting group is removable by catalytic hydrogenation, then that protecting group will be removed by the hydrogenation conditions of the reaction of Scheme 3. That removal, forming the compounds of formula V(b) where R4 is hydroxy, facilitates the installation of the R4 group found in the antifolate final products, i.e., the chiral glutamic acid, discussed below.\nCompounds of formula V which possess antifolate activity are those where R4 is NHC*H(C(O)R5)CH2CH2C(O)R5) and R6 is hydrogen (hereafter referred to as xe2x80x9cfinal productsxe2x80x9d), and thus, those compounds are preferred. Although the glutamate side chain can be installed at any point in the overall process of this invention, when the processes of Schemes 1, 2, 3 are performed in sequence, with the preferred groups noted above, a preffered time to install the glutamate residue is after performing the hydrogenation described in Scheme 3. This is accomplished by coupling a compound of formula V or V(b) where R4 is hydroxy with a carboxy protected glutamic acid derivative of the formula H2NC*H(C(O)R5xe2x80x2)CH2CH2C(O)R5xe2x80x2, in the manner generally described in PCT application WO 86/05181, utilizing conventional condensation techniques for forming peptide bonds. These techniques include activation of the carboxy group through formation of a mixed anhydride or acid chloride, treatment with dicyclohexylcarbodiimide, or use of diphenylchlorophosponate. For further instruction on general methods of forming this amide bond, see, e.g., Bodanszky, M., Principles of Peptide Synthesis, 2nd Ed., Springer-Verlag, Berlin, Heidelberg, 1993. It is preferred that the glutamate side chain is present and that the R5 groups found in that side chain are both carboxy protecting groups, e.g., C1-C6 alkoxy, when performing the processes of Schemes 4 and 5. All discussions and structures pertaining to Schemes 4 and 5 below relate to the situation where the preferred substituents at R4 are present but those substituents are not required for the processes of Scheme 4 and 5 to be operable.\nOf the final product compounds of formula V, the compounds of formula V(c)(R), shown below, are preferred due to enhanced antifolate activity relative to the compounds of formula V(c)(S). These individual diastereomeric final products, prepared as described above or in the previously incorporated by reference U.S. Pat. Nos. 4,684,653 and 4,882,334, may be separated as taught in those patents, i.e., by chromatography or preferably recrystallization. For example, an appropriately selected chiral acid may be employed to form a mixture of diastereomeric salts more amenable to selective recrystallization of one diastereomer as illustrated in Scheme 4 below where R3, R5xe2x80x2, R6, R7, and Y are as defined above. \nIn order to carry out the separation of Scheme 4, it is necessary that R6 in compounds of formula V(c) is hydrogen. If the process of Scheme 4 is performed with the preferred group at R6, i.e., a protecting group removable by acid hydrolysis such as 2,2-dimetyl-1-oxopropyl, then that protecting group will be removed by the addition of the chiral acid and a separate step to remove it will not be necessary. Thus, compounds of formula V(c) where R6 is hydrogen do not necessarily have to be prepared in a separate step before proceeding to the resolution process of Scheme 4. The elimination of the requirement for this extra step is why acid labile amino protecting groups are preferred.\nThe first step in Scheme 4 is the addition of a chiral acid to a mixture of compounds of formula V(c) dissolved or suspended in a suitable solvent. This addition performs two functions: the chiral acid removes the amino protecting group at R6 and forms a diastereomeric acid addition salt of the compound of formula V(d). When Y is CHxe2x95x90CH, a preferred acid for this purpose is (1S)-(+)-camphorsulfonic acid. When Y is O or S, a preferred acid is (1R)-(xe2x88x92)-camphorsulfonic acid. A preferred solvent for the removal of the protecting group and formation of the salt is a lower alcohol preferably ethanol.\nOnce the salt is formed, the separation or recrystallization of Step 2 is performed by suspending the compounds of formula V(d) in a suitable solvent, heating the mixture until a solution is formed, and then allowing the solution to cool in order to precipitate the desired isomer. The important parameters in a chiral resolution in general, and when specifically resolving the compounds of formula V(d), are the solvent system, stir rate, and temperature. Preferred solvent systems are mixtures of a lower alcohol, preferably ethanol, and water. The ratio of ethanol to water by volume can be from about 0.33 to about 3 to 1, but a 1:1 mixture is preferred. The ratio of solvent to solute should be about 10 to about 20 to 1 but the preferred ratio is about 15 to 1. The rate of stirring during crystallization can have a marked effect on the resolution. It is preferred that once the salts are formed and dissolved by heating in the crystallization solvent, that the samples are not stirred while cooling. Temperature may also have an impact on resolution. Continued cooling below ambient temperature can increase the recovery of product but at the expense of separation efficiency. It is preferred to allow the crystallization to occur at a temperature between about 20xc2x0 C. and 34xc2x0 C. with ambient temperature (about 22xc2x0 C.) being most preferred.\nIn order to avoid hydrolysis of the esters on the glutamate residue by the aqueous acidic conditions of the resolution, a buffer such as sodium acetate is preferably employed. An amount of sodium acetate approximately equal to the excess of the acid is preferably added after solvolysis of the amino protecting group is complete.\nThe filtrate produced in Scheme 4 will contain a mixture of diastereomeric salts of formula V(d). This mixture is enhanced with the diastereomer of formula V(d)(S), i.e., the acid addition salt of V(d)(S), that didn\"\"t crystallize. This enhanced mixture is not usually amenable to further separation by crystallization. Thus, a large majority of the desired isomer was heretofore unrecoverable by further crystallization. Scheme 5 below, where R3, R4xe2x80x2, R6, R7, and Y are as defined above, illustrates another novel method of preparing compounds of formula IV which facilitates the further separation of this mixture by crystallization. \nIf the process of Scheme 5 is to be performed using the contents of the filtrate from Scheme 4 as a source of starting materials, it is preferred that the free base compounds of formula V(c) are first extracted from that mixture. This is easily accomplished by diluting the mixture with a water immiscible organic solvent and treating the mixture with a weak base dissolved in water such as a solution of sodium bicarbonate in water. Dichloromethane and aqueous sodium bicarbonate are preferred reagents for this purpose. For example, treatment of the filtrate with dichloromethane (about 0.33 mL/g filtrate) and about 2 equivalents of 1M aqueous sodium bicarbonate for about 30 minutes at room temperature will afford two clear, readily-separable phases. The organic phase is separated and preferably dried with a common drying agent or agents before proceeding. Sequential washes with brine and 1M aqueous sodium bicarbonate is a preferred drying procedure.\nOnce the extraction procedure is performed, an amino protecting group must be reinstalled at R6 or that amino group must be protonated before performing the process illustrated in Scheme 5. Reprotection is preferred and may be accomplished as taught in Greene cited above or as discussed in Preparation 9 below. It is preferred that the amino protecting group be removable by acid hydrolysis, e.g., a 2,2-dimethyl-1-oxopropyl protecting group. The resulting compounds of formula V(e) are then separated as follows.\nSingle diastereomers or mixtures of any ratio of compounds of formula V(e) may be dissolved or suspended in a suitable solvent and an oxidizing reagent added to provide the compounds of formula IV(b). Choice of solvent, reaction temperatures, and times will depend generally on the oxidizing reagent employed.\nDDQ is a preferred oxidizing agent. When DDQ or chloroanil is employed as the oxidizing agent, hydrocarbons such as pentane, hexane, toluene, and the like; lower alcohols such as methanol, ethanol, isopropanol, and the like; or chlorinated hydrocarbons such as chloroform, dichloromethane, and the like; are suitable. Chlorinated hydrocarbons, especially dichloromethane, are preferred. When the oxidation is performed under the preferred conditions on the preferred compounds of formula V(e), chromatography to purify the resulting compounds of formula IV(b) is generally not necessary. See, e.g., Example 6 below.\nThe reaction is typically allowed to proceed at temperatures between 0xc2x0 C. and 200xc2x0 C. for from about 30 minutes to about 24 hours. The reaction is preferably performed at temperatures between 15xc2x0 C. to 80xc2x0 C. Even more preferred is when the reaction is performed between 20xc2x0 C. to 40xc2x0 C., and most preferred is when the reaction is performed at room temperature for from 20 minutes to 1 hour.\nThe amount of oxidizing reagent will vary depending on which oxidant is employed but will generally range from about 0.1 equivalents to about 5 equivalents relative to the compound of formula V(e). When DDQ is employed, about 1.1 to about 3 equivalents are preferred. Even more preferred is from 1.8 to about 2.2 equivalents while 1.9 to about 2.1 equivalents is most preferred.\nThe carboxy and/or amino protecting groups in compounds of formula IV(b) are preferably not removed as taught in Greene. Instead, a preferred course of action is to reduce the compounds of formula IV(b) back to a 50:50 diastereomeric mixture of compounds of formula V(e) in order to perform the chiral acid separation taught in Scheme 4.\nThe hydrogenation may be performed by dissolving or suspending a compound of formula IV(b) in a suitable solvent, in the presence of a hydrogenation catalyst, and exposing the mixture to an atmosphere of hydrogen. A convenient and preferred solvent is about a 4:1 mixture by volume of tetrahydrofuran and ethanol. A convenient and preferred hydrogenation catalyst is 5% palladium on carbon. The catalyst is typically employed, relative by weight to compounds of formula IV(b), in a range of from about 10% to about 200%. Preferably, the range is from about 20% to about 75%, with about 25% being preferred. It is typically preferred to create an atmosphere of hydrogen where the pressure of hydrogen is equal to or greater than that of ambient pressure. A typical pressure range is from about ambient to about 100 psi of hydrogen. More preferred is an atmosphere of hydrogen between 40 psi and about 60 psi with 50 psi most preferred. The reaction may be performed at temperatures ranging from ambient to about the reflux temperature of the mixture. At 50 psi of hydrogen, the preferred reaction temperature is about 100xc2x0 C., with the reaction typically substantially complete in less than 4 hours.\nOnce a 1:1 mixture of compounds of formula V(e) is formed as described in Scheme 5, it is further amenable to separation as discussed in Scheme 4 above. This cycle may be repeated as many times as the practitioner wishes in order to maximize the yield of the diastereomer of formula V(d)(R).\n2,4-Diamino-6-hydroxypyrimidine, protected glutamic acids of the formula H2NC*H(C(O)R5xe2x80x2)CH2CH2C(O)R5xe2x80x2 and compounds of formula I and II are known in the art, and, to the extent not commercially available are readily synthesized by standard procedures commonly employed in the art. For example, see Preparations 1-6 below.\nThe optimal time for performing the reactions of Schemes 1-5 can be determined by monitoring the progress of the reaction by conventional chromatographic techniques. Choice of solvent is generally not critical so long as the solvent employed is inert to the ongoing reaction and sufficiently solubilizes the reactants to afford a medium within which to effect the desired reaction. Unless otherwise indicated, all of the reactions described herein are preferably conducted under an inert atmosphere. A preferred inert atmosphere is nitrogen."}
{"text": "Growing demand for media content means there is a growing need for bandwidth not just to, but within, a customer premises. For example, for playback of high resolution video on wireless devices, needed bandwidth, e.g., via existing radio-based technologies such as LTE (long term evolution) wireless communications, and Wi-Fi communications, i.e., according to IEEE 802.11, is often not available. Another technology for supporting high bandwidth media transmissions is Li-Fi, i.e., Light Fidelity, communications. Using the visible light spectrum, Li-Fi technology can transmit data thousands of times faster than existing radio spectrum based technologies. However, Li-Fi requires that a receiving Li-Fi device receives light beams transmitted by a Li-Fi transmitter, i.e., photo detectors of a Li-Fi receiver have to be able to actually “see” the transmitted light in order to capture the data."}
{"text": "As a next-generation nonvolatile memory substituting for a flash memory or the like which has come to have a limit in miniaturization, a resistance change memory for storing data by using a resistance variable element, such as MRAM (Magnetoresistive Random Access Memory), ReRAM (Resistance Random Access Memory), PCRAM (Phase Change Random Access Memory), and the like has attracted attention.\nAs a method of increasing the memory density (increasing capacity), there is a method of reducing the size of the elements constituting the memory, as well as a method of multileveling the recording bits per element constituting the memory, and various multileveling methods have been proposed (for example, see Patent Documents 1 to 3).\nThere is one type of MRAM called a magnetic wall driving type or a magnetic wall shifting type (for example, see Patent Document 4). The magnetic wall driving MRAM allows current to flow in an in-plane direction of a magnetic wall driving layer (or magnetization free layer), shifts the magnetic wall by the spin transfer effect of spin polarized electrons, and reverses the magnetization direction of the ferromagnetic film according to the direction of the write current and writes data.\nPatent Document 4 describes a method of multilevel recording and analog recording with respect to a magnetic wall driving MRAM.\nIn MRAM, different writing methods of data have been proposed, and in addition to a magnetic wall driving MRAM, a magnetic field writing MRAM, a yoke magnetic field writing MRAM, an STT (Spin Transfer Torque) MRAM, an SOT (Spin Orbit Torque) MRAM, and the like are known."}
{"text": "Automated production of assembly parts, for example, requires some means for identifying, inspecting, and measuring the parts at high speed. The use of a video camera to scan each part visually together with a digital system to analyze the video signal generated by the camera is well-known. One known prior art system converts each image pixel to a \"black\" or \"white\" binary coded bit. During a frame scan, this stream of bits is gated by a window generator to an accumulator. The window generator gates only the bits corresponding to pixels within a preselected window area within the camera field. The size of the window can also be preselected. The accumulator counts all the bits within the window that are of one value (either black or white). By providing multiple windows generators and at least as many accumulators, the black or white pixel count from more than one window can be counted during one frame scan. However, the number of windows is limited by the number of accumulators, and an accumulator cannot be assigned to more than one window at a time. Such a system is severely limited in the number and size of windows that can be used in the system by the number and capacity of the accumulators."}
{"text": "1. Field of the Invention\nThe present invention relates to a digital video processor, and more particularly to a secondary digital color processor (DCP) which is usable with modern standard-definition digital telecines and is also compatible with high-definition television systems, for correcting video color and for other purposes.\n2. Background Art\nTelecines and external secondary color correctors (ESCC's) are used for the production of high-quality television pictures in a post-production television studio environment. The highest quality pictures normally originate from motion picture film, normally 35 mm film. These pictures are transcribed to video on a telecine, such as the Rank Cintel MkIII, the Rank Cintel URSA, or the BTS FDL90. Telecine machines convert the picture to a digital format and contain processing elements to enable simple manipulation of the color reproduction of the resulting images.\nThis sort of color manipulation may also be performed by an external analog system such as the RCA Chromacomp or its derivatives. This simple RCA-type color processing has the limited capability of, for example, intensifying the red component, over the entire picture, in all colors that have a red component. However, it is incapable of, for example, altering only the dark reds and leaving the light reds unchanged.\nOver the years, improved ESCC's such as the PRISM from Encore Video Industries, Inc., of California, have emerged. The PRISM is an analog device which is capable of more selective color adjustment than the RCA Chromacomp. Other analog ESCC's are the Da Vinci, and the \"Sunburst\" made by Corporate Communications Consultants.\nColor correctors have also been used for many years for tape-to-tape correction. The output of one videotape recorder (VTR) is processed by the ESCC, whose output is then recorded by another VTR.\nESCC's were traditionally analog, but now some color correctors are available that are partly or fully digital. The output of most modern telecines is digital.\nAnalog color correctors have several problems. Since a logical way to interconnect the telecine to the color corrector is to connect the output of the telecine to the input of the color corrector, it has normally been necessary to convert the digital signal to analog, apply the color correction, and then convert the corrected signal back to digital for storage or further processing. The conversion from digital to analog and back again causes degradation, and is therefore undesirable.\nThis phenomenon is most easily visible when one feeds a color signal through a color corrector with all of the controls set to zero (which in theory, should not change the colors). When one compares the output of the color corrector with a signal which bypasses the color corrector, quite often, there will be a difference between these two signals, which will become apparent when the two images are viewed \"split-screen\".\nOne reason for this discrepancy is that the signals in the electronic signal path (particularly in analog) are unintentionally modified due to the imperfections of real electronic circuits.\nAnother reason is that since the color corrector is an analog device, it will suffer from drift, which is inevitable in analog devices. Drift originates from many sources, one of which is temperature.\nA third reason for the visible differences is noise. This noise may be visible as a change in level (and therefore color) or a difference in texture of a given flat color bar.\nEven if, as has occasionally been done, the signals are processed in a totally digital domain, there have been degradations.\nOne reason may be digital \"rounding\", such that a color passing through a color corrector unintentionally becomes slightly modified, compared with the input.\nFurther, in known color correction systems (particularly digital systems), discontinuities occur. This happens where colors below a given limit are to remain unmodified, but colors above that limit are to be modified. If the picture to be modified contains an increasing ramp of that color, then the portion of the picture around the limit will show an unnatural discontinuity when it is color-corrected.\nThe known color-correction devices will be discussed further below, in the context of the following discussion of their use with telecine devices.\nReferring first to FIG. 1, there is seen a schematic block diagram of the most basic type of stand-alone telecine 20, which is usable to broadcast a film directly over the air in real time. A film transport 21 runs a film 22 past a CRT scanner 24, which incorporates photomultipliers for generating red, green, and blue component signals designated R.sub.v, G.sub.v, and B.sub.v. Gain, lift and gamma of the three signals are adjusted by a processor 26, and color is adjusted by a processor 28. Simple overall picture adjustments such as lightening and darkening are provided by a local control panel 30. Adjustments are performed live, as the film is broadcast. The local control panel 30 also contains controls such as a START control for the film transport and an ON control for the CRT scanner. A disadvantage of this basic system having only the foregoing components is that, in order to provide interlace for producing a conventional broadcast signal, it is necessary for the lines in each film frame to be scanned non-sequentially, that is, lines 1, 3, 5, . . . , followed by lines 2, 4, 6 . . .\nAn improvement upon the foregoing basic system, developed in the late 1970's, is the addition of a store 32, also shown in FIG. 1. The store improves the scanning process, by permitting the lines of each frame to be scanned sequentially, and then read out of the store with interlace. This feature was patented by Rank Cintel in British patents 1,542,213 and 1,535,563, and equivalent U.S. Pat. Nos. 4,184,177 and 4,127,869, respectively.\nAn improvement shown in FIG. 2 is a controller/programmer 30', which as shown is an accessory for the telecine 20', although it could conceivably be built into the telecine 20'. The controller/programmer 30' provides a control panel for the processors 26, 28, and the film transport 21. It further includes a programming function, whereby an operator can rehearse, slowly, the optimum grading and picture adjustments for each scene in the film, which the programmer stores. A key feature is that the programmer is driven by time code, or by film footage, so it can tell one scene from another by the time code or by the position of the film.\nOne example of a controller/programmer is the POGLE telecine controller/programmer manufactured by Pandora International Ltd., which is designed for use with a range of telecine machines. It is capable of providing a large number of control signals simultaneously, for example 96 control channels, either analog (a voltage within a predeterminable range) or digital. The channels can be used to control a telecine and/or other peripheral devices, such as noise reducers, VTR's, still stores, mixers, and color correctors (such as the DCP disclosed herein).\nAnother example of a controller/programmer is the \"Rainbow\" system of Corporate Communications Consultants.\nA further improvement, shown in FIG. 3, is the external secondary color corrector (ESCC) 34. Examples are the RCA Chromacomp, the PRISM, and the Da Vinci. ESCC's can be either digital or analog. Instead of simple RGB control, which only gives the ability to add or subtract red, green, and blue, everywhere in the picture, this advanced generation of ESCC's provides 6-or-more-channel color control. That is, for example, 6 separate colors can be selected and then modified, without modifying any other colors. The ESCC 34 is controlled by the controller/programmer 30\".\nThe ESCC 34 takes its inputs from the color processor 28 in the telecine 20\", and provides its outputs to the store 32. The reason for this signal path is that the store is usually capable of handling lower bandwidth (i.e., fewer bits or lower resolution) than the signals from the processor 28, so the number of bits is reduced, for example, by rounding, before the signals are stored. For example, the URSA's output is medium-bandwidth \"4:2:2\" color (D-1 format). Thus, taking the input to the ESCC 34 from the output of the processor 28 allows the ESCC to operate on a higher-bandwidth signal than if it operated on the output from the store.\nThe trend is for telecines to be made digital and to incorporate improved color correction facilities. For example, one advanced digital telecine, the Rank Cintel URSA telecine, has a built-in digital implementation of a 6-vector RCA-type processor, in place of the analog processor 28 shown in FIGS. 1-3. It would be desirable to provide an ESCC which is capable of processing the digital output from the internal digital color processor 28 in the URSA (or the like), having greater capacity and flexibility than that RCA-type color corrector.\nA further trend is to increase the bandwidth of the store 32, which will make it practical for the digital color processor to perform its functions on the digital output of the store, making it unnecessary to connect the ESCC to any internal circuits of the telecine.\nIn either case, an ESCC is needed which will not noticeably reduce the bandwidth of the digital color signal or degrade the picture.\nThe disclosures of all prior art publications mentioned herein are expressly incorporated by reference."}
{"text": "Radio-Frequency Identification (RFID) systems typically include RFID readers, also known as RFID reader/writers or RFID interrogators, and RFID tags. RFID systems can be used in many ways for locating and identifying objects to which the tags are attached. RFID systems are particularly useful in product-related and service-related industries for tracking objects being processed, inventoried, or handled. In such cases, an RFID tag is usually attached to an individual item, or to its package.\nRFID techniques entail using an RFID reader to interrogate one or more RFID tags. The reader transmitting a Radio Frequency (RF) wave performs the interrogation. The RF wave is typically electromagnetic, at least in the far field. The RF wave can also be predominantly electric or magnetic in the near field. The RF wave may encode one or more commands that instruct the tags to perform one or more actions.\nA tag that senses the interrogating RF wave responds by transmitting back another RF wave. The tag generates the transmitted back RF wave either originally, or by reflecting back a portion of the interrogating RF wave in a process known as backscatter. Backscatter may take place in a number of ways.\nThe reflected-back RF wave may encode data stored in the tag, such as a number. The response is demodulated and decoded by the reader, which thereby identifies, counts, or otherwise interacts with the associated item. The decoded data can denote a serial number, a price, a date, a time, a destination, an encrypted message, an electronic signature, other attribute(s), any combination of attributes, and so on. The data items are sometimes known as codes. Accordingly, when a reader receives tag data (codes) it can learn about the item that hosts the tag and/or about the tag itself.\nAn RFID tag typically includes an antenna system, a radio section, a power management section, and frequently a logical section, a memory, or both. In some RFID tags the logical section may include a cryptographic algorithm which, in many instances, relies on one or more passwords or keys stored in tag memory. In earlier RFID tags the power management section included an energy storage device such as a battery. RFID tags with an energy storage device are known as battery-assisted, semi-active, or active tags. Advances in semiconductor technology have miniaturized the electronics so much that an RFID tag can be powered solely by the RF signal it receives. Such RFID tags do not include an energy storage device and are called passive tags. Of course, even passive tags typically include temporary energy- and data/flag-storage elements such as capacitors or inductors.\nA well-known problem in RFID systems is speed in reading the tags, especially when the reader reads more than one code from each tag. The problem is exacerbated if there are many tags, or if the host items are moving and thus allow only limited time to read their tags."}
{"text": "Although the Actinomycetales produce more than half of the known antibiotics having valuable clinical and other applications as secondary metabolites and, thus, are recognized as a key target for application of gene manipulation techniques, many problems remain to be overcome before specific useful genes are successfully identified and cloned [\"Molecular Breeding and Genetics of Applied Microorganisms\", Sakaguchi and Okanishi, eds., Academic Press (New York) Kodansha Ltd. (Toyko) 1980, pgs. 130-131]. Until the present work, cloning of a .beta.-galactosidase gene from a Streptomyces species onto a suitable vector followed by introduction and expression of such vector has not been reported. Prior work has concerned development of other cloning systems or vectors for Streptomycetes [Bibb et al. (1978), Nature 274: 398-400; Hayakawa et al. (1979), J. Antibiot. XXXII(12): 1348-1350; Okanishi et al. (1980), J. Antibiot. XXXIII(1): 88-91; Bibb et al. (1980), Nature 284: 526-531; Thompson et al. (1980), Nature 286: 525-527; Suarez et al. (1980), Nature 286: 527-529; Bibb et al. (1981), Mol. Gen. Genet. 184: 230-240]; [Bibb (1981), \"Microbiology-1981\", Schlessinger, ed., American Society for Microbiology, (Washington, D.C.) 1981, pgs. 367-370 and Hopwood et al. (1981), \"Microbiology-1981\", supra. pgs. 376-379], cloning and expression in Streptomyces sp. of genes derived from Escherichia coli [Schottel et al. (1981), J. Bacteriol. 146: 360-368] and cloning of genes from Streptomycetes in Escherichia coli [\"Molecular Breeding and Genetics of Applied Microorganisms\", supra; pgs. 130-137]. Chater et al. (1982), Current Topics in Microbiol. and Immunol. 96: 69-95, review gene cloning in Streptomyces and is incorporated by reference herein as though fully set forth.\nWork with various .beta.-galactosidase genes, their expression and application of such expression as an assay or detection method has been reported by Rose et al. (1981), Proc. Natl. Acad. Sci. USA 78(4): 2460-2464, for expression in yeast of yeast genes fused to .beta.-galactosidase genes from Escherichia coli; by Casadaban et al. (1980), J. Mol. Biol. 138: 179-207, for fusion of .beta.-galactosidase genes to promoters in Escherichia coli and assay following transformation; and by Talmadge et al. (1981), Nature 294: 176-178, for construction of Escherichia coli containing a plasmid encoding a .beta.-galactosidase-preproinsulin fusion protein.\nCollinge et al., U.S. Pat. No. 3,816,259, disclose that a Streptomyces coelicolor preparation had a .beta.-galactosidase activity.\nIn general, the activity of promoters can be assayed by measuring the amount of gene product which is formed as a consequence of transcription starting from a specific promoter. The amount of gene product formed is determined by using a specific property of that gene product, such as enzymatic activity. In studying gene expression or in constructing high expression vectors which rely on highly efficient promoters, the gene which is naturally expressed from such a promoter is replaced by the structural gene whose product can be more easily monitored. The lacZ gene from Escherichia coli [Casadaban et al. (1980), supra.] is frequently used for this purpose.\nA variety of chromogenic substrates, such as 5-bromo-4-chloro-3-indolyl-.beta.-D-galactosidase (referred to as \"X-gal\") or o-nitrophenyl-.beta.-D-galactoside (referred to as \"ONPG\") can be used to monitor enzymatic activity as described by Miller (1972), \"Experiments in Molecular Genetics\", Cold Spring Harbor Laboratories (Cold Spring Harbor, N.Y.). These substrates are advantageous since the efficiency of a promoter fused to a gene coding for an enzyme which can react with the substrate, such as the lacZ gene, can be monitored by growing the organism on a solid agar medium containing the substrate and observing for enzyme-substrate reaction. In this manner, several hundred individual colonies can be scored at one time for their ability to express the gene. Thus, relatively rare events such as the occurrence of a highly efficient promoter can be detected. .beta.-Galactosidase expression can be used in such a procedure to assay gene transcription and to detect and isolate mutants which over-produce a particularly desired protein, such as an enzyme involved in antibiotic production.\nIn order to effectively use such a powerful approach as described above, it is crucial that the chosen substrate has the opportunity to react with the enzyme. If, as in the case of .beta.-galactosidase produced by Escherichia coli, the enzyme is intracellular, the substrate must enter the cell in order for the enzyme-substrate reaction to occur. With Streptomyces lividans, however, the commonly used substrates, X-gal and ONPG, enter the cell only poorly as verified by comparing the intracellular .beta.-galactosidase activity of a suitable organism with dye formation on plates. For example, we have found that although intracellular activity as measured with cell extracts and with ONPG as substrate was very high (300 nmoles/mg protein/min), no significant color reaction with whole cells and with either ONPG or X-gal was found. Furthermore, Actinomycetes differ from many other microorganisms by the formation of an aerial mycelium which separates the cells physically from the substrate, thus further restricting access of the substrate to the cells [Kalakoutskii et al. (1976), Bacteriol. Rev. 40(2): 469-524]."}
{"text": "This invention concerns certain derivatives of fluorenes that are useful as antiarrhythmic agents.\nU.S. Pat. Nos. 4,197,313 and 4,277,495 describe various 9-aminoalkylfluorene derivatives that are useful as antiarrhythmic agents. All of the disclosed compounds are primary, secondary or tertiary amines, or acid addition salts thereof. It has long been known that many antiarrhythmic agents also have local anesthetic activity, and that the cardiac activity of such compounds may in some way be associated with the local anesthetic properties. It has been thought that in order to have anesthetic and/or antiarrhythmic activity, an amino compound had to be basic, ie. a primary, secondary or tertiary amine. Accordingly, no non-basic aminoalkyl fluorene antiarrhythmic agents have heretofore been known.\nAn object of this invention is to provide a group of non-basic N-alkanoyl-aminoalkylfluorenes that are active as antiarrhythmic agents."}
{"text": "This specification relates to determining whether a submission of data by a user is accurate.\nA search system can provide one or more knowledge panels in response to a received search query. A knowledge panel is a user interface element that provides a collection of information or other content related to a particular entity referenced by the search query. For example, the entity may be a person, place, country, landmark, animal, historical event, organization, business, sports team, sporting event, movie, song, album, game, work of art, or any other entity.\nIn general, a knowledge panel provides a summary of information about the entity. For example, a knowledge panel for a famous singer may include the name of the singer, an image of the singer, a description of the singer, one or more facts about the singer, content that identifies songs and albums recorded by the singer, and/or links to searches related to the singer. Other types of information and content can also be presented in the knowledge panel. Information presented in a knowledge panel can include content obtained from multiple disparate sources, e.g., multiple different web pages accessible over the Internet.\nA search system can maintain a knowledge base that stores information about various entities. The system can assign a unique entity identifier to each entity. The system can also assign one or more text string aliases to a particular entity. For example, the Statue of Liberty can be associated with aliases “the Statue of Liberty” and “Lady Liberty.” Aliases need not be unique among entities. For example, “jaguar” can be an alias both for an animal and for a car manufacturer.\nThe system can also store information about an entity's relationship to other entities. For example, the system can define a “located in:” relationship between two entities to reflect, for example, that the Statue of Liberty is located in New York City. In some implementations, the system stores relationships between entities in a representation of a graph in which nodes represent distinct entities and links between nodes represent relationships between the entities. In this example, the system could maintain a node representing the Statue of Liberty, a node representing New York City, and a link between the nodes to represent that the Statue of Liberty is located in New York City."}
{"text": "This invention generally relates to explorations for hydrocarbons involving electrical investigations of a borehole penetrating an earth formation. More specifically, this invention relates to performing such borehole investigations employing a logging fluid which is an invert emulsion or a non-aqueous fluid, wherein the logging fluid has an electrical stability of at most 350 volts as measure by API RP 13B-2. In one aspect of the present invention, this invention relates to performing such borehole investigations employing a logging fluid comprising a non-conductive fluid which has been modified by the addition of an additive, wherein the additive is a salt having a cation, preferably a divalent or trivalent cation, which is capable of forming a complex, in an effective amount to sufficiently modify the electrical characteristics thereof resulting in a modified fluid which when disposed in at least a portion of the borehole can be used with a conventional electrical well logging tool to perform such investigations as the tool is moved along such portion of the borehole.\nDuring or after the drilling of oil or gas wells, measurements of the electrical characteristics of the wellbore are performed. Electrical earth borehole logging is well known and various devices and techniques have been described. A variety of measurements may be made, but typically include resistive measurements extending deep into the formation and also superficial measurements of changes in resistivity at the surface of the borehole. In an electrical investigation of a borehole, current from an electrode is introduced in the formation from a tool inside the borehole. If this current is maintained constant, the voltage measured at a monitor electrode is proportional to the resistivity of the earth formation being investigated. If the current is varied to maintain constant the voltage measured at a monitor electrode, the current is inversely proportional to the resistivity of the earth formation being investigated. If both voltage and current are allowed to vary, their ratio is proportional to the resistivity of the earth formation being investigated. Substantial advances have been made in such electrical investigations by using electrodes whose currents are focused by other electrodes and thus determine the resistivity of the formation at a desired distance from the borehole wall surface. Examples of such techniques and devices for focused electrical investigation are described and shown in the U.S. Pat. Nos. 2,712,629 to Doll; 2,750,557 to Bricuad; 3,521,154 to Maricelli; and 4,468,623 to Gianzero et al.\nIn U.S. Pat. No. 2,712,629 to Doll, pad mounted sets of electrodes are described as each formed of a central survey electrode surrounded at spaced intervals by continuous guard electrodes embedded in segmented recesses.\nIn U.S. Pat. No. 2,750,557 to Bricuad, the pad mounted electrodes are formed of electrically directly connected segments or buttons.\nIn U.S. Pat. No. 3,521,154 to Maricelli, a plurality of survey electrodes are mounted on a single pad as a composite focusing electrode, with a pair of the survey electrodes aligned along the direction of travel of the tool along the borehole and one survey electrode horizontally displaced to provide a technique for effectively improving the signal to noise ratio of the resistivity measurements.\nIn U.S. Pat. No. 4,468,623 to Gianzero et al., an earth formation investigating tool is described in which borehole wall features on the order of millimeters in size can be detected. The tool includes an array of small cross-section survey electrodes (buttons ) which are pressed towards the borehole wall and each button injects an electric current into the adjoining formation. The individual button currents are monitored and signals representative of button currents are recorded as curves as a function of depth. The measured button currents reflect the resistivity of the material in front of each button. In order to achieve a high resolution investigation, the electrodes are arranged in an array of multiple rows. The electrodes are so placed at intervals along a circumferential direction about the borehole axis as to inject survey currents into borehole wall segments which overlap with each other to a predetermined extent as the tool is moved along the borehole wall. In this manner, a detail high resolution resistivity or conductivity investigation of the borehole wall can be made. The presence of a fracture may be identified by noting a deviation between the survey currents from different pads. Such survey current deviation may indicate a fracture by virtue, for example, of the invasion of higher conducting mud into the fracture.\nAs inferred from the foregoing reference to xe2x80x9chigher conducting mudxe2x80x9d, such logging tools primarily were designed for use in an electrically conductive aqueous-based fluid. Accordingly, it has been possible to obtain electric logs from well boreholes primarily by suspending the logging tool in an electrically conductive aqueous media. As oil wells were and continue to be drilled deeper into water sensitive formations, the application of non-aqueous-based drilling fluids or invert emulsions of water or brine in the various types of fluids used in such non-aqueous-based drilling fluids, such as organic solvents, diesel fuel, mineral oil, vegetable oil and synthetic fluids, is increasing. Because the continuous phase of this invert emulsion is usually a non-conductor of electricity, conventional electric logs which require passage of electrons through a conductive media have not been effective in these types of drilling and completion fluids.\nSome logging tools have been designed for use in oil-based drilling fluids. For example, in such tools, knife-edge electrodes have been used to ensure contact with the filter-cake or mudcake which usually forms on the side of the borehole. However, U.S. Pat. No. 3,521,154 to Maricelli notes that since the oil-based drilling fluid or mud is relatively non-conductive, even the slightest separation between the knife-edge and the mudcake will provide erroneous indications of the conductivity of the adjoining formation. In U.S. Pat. No. 2,930,969 to Baker, the tool thereof may employ brush-like contacts which scratch through the mudcake and effect good electrical connection with the rock when an oil-based fluid is used.\nA few attempts to make oil-based drilling fluids electrically conductive for the purpose of electrical logging have been reported though none of them has been a commercial success. U.S. Pat. No. 2,696,468 to Fischer disclosed a conductive oil-based drilling fluid containing up to 10 percent by weight water, an electrolyte and certain types of emulsifying agents, specifically sulfated and sulfonated organic compounds which promote the formation of oil-in-water emulsions. The electrolytes were water-soluble ionizable metallic compounds and were for the most part water-soluble salts of alkali- and alkaline-earth metals and alkali-metal hydroxides. Though almost forty compounds, including magnesium chloride, magnesium nitrate and magnesium sulfate, were specifically named as electrolytes, only seven sodium-containing compounds and calcium chloride were exemplified. Fischer disclosed a particular preference to alkali-metal hydroxides, silicates and phosphates. U.S. Pat. No. 2,739,120 also to Fischer discloses similar oil-based fluids which are asserted to be electrically conductive and contains a non-ionic surfactant rather than the emulsifiers of U.S. Pat. No. 2,696,468. Though almost forty compounds, including magnesium chloride, magnesium nitrate and magnesium sulfate, are specifically named as electrolytes, only sodium-containing compounds were exemplified. Both of these patents disclosed that in order to maintain the general desirable characteristics of oil-based drilling fluids, the water content should be maintained below 10 percent by weight, i.e., avoiding the formation of a water-in-oil or invert emulsion.\nAbout twenty-five years later, Hayes et al. in U.S. Pat. No. 4,012,329 disclosed water-in-oil microemulsion drilling fluids which were asserted as being capable of conducting electrical current and as such permitted the use of ordinary electrical logging techniques. This fluid contained water, sodium petroleum sulfonate, hydrocarbon, bentonite and, optionally, cosurfactant, electrolyte, gelling agents and fluid loss agents. The electrolyte was a water-soluble inorganic base, inorganic acid or, preferably, inorganic salt. Certain sodium and potassium salts were identified, but not specifically identified in the examples.\nTherefore, a need exists to modify non-conducting fluids in order to effectively use the great variety of conventional electrical well logging tools, particularly imaging tools.\nAccordingly, a well logging fluid comprising:\na relatively non-conductive fluid, wherein the non-conductive fluid is selected from the group consisting of a non-aqueous fluid and an invert emulsion of an aqueous phase in a non-conductive liquid, wherein the aqueous portion of the invert emulsion ranges up to about 70 percent by volume, preferably from about 10 to about 40 percent by volume; and\nan effective amount of a salt sufficient to modify the electrical characteristics the relatively non-conductive fluid resulting in a modified fluid or well logging fluid such that an electrical well-logging tool when disposed in the modified fluid can perform an electrical well-logging survey of a borehole containing the modified fluid,\nwherein the salt has a cation and an anion, the cation is selected from the group consisting of divalent and trivalent metal cations capable of forming a complex, and the anion is one that disassociates from the cation in an aqueous medium. Preferably, the modified fluid or well logging fluid has an electrical stability of at most 350 volts, more preferably at most 250 volts, as measured according to API RB 13-2. The metals for the cation include magnesium, zinc, cobalt, copper, aluminum and iron, preferably magnesium, zinc and aluminum. The anion includes halides, nitrates, formates, citrates and acetates, preferably halides and nitrates.\nThere is also provided a method for measuring a characteristic of earth formations traversed by a borehole, the method comprising the steps of:\nproviding a relatively non-conductive fluid, wherein the non-conductive fluid is selected from the group consisting of a non-aqueous fluid and an invert emulsion of an aqueous phase in a non-conductive liquid, wherein the aqueous portion of the invert emulsion ranges up to about 70 percent by volume, preferably from about 10 to about 40 percent by volume;\nadding an effective amount of a salt to the non-conductive fluid sufficient to modify the electrical characteristics thereof resulting in a modified fluid such that an electrical well-logging tool when disposed in the modified fluid can perform an electrical well-logging survey of a borehole containing the modified fluid,\nwherein the salt has a cation and an anion, the cation is selected from the group consisting of divalent and trivalent metal cations capable of forming a complex, and the anion is one that disassociates from the cation in an aqueous medium;\nproviding a column of the modified fluid in a portion of the borehole;\ndisposing an electrical well-logging survey tool in the portion of the borehole; and\nperforming an electrical well-logging survey of the portion of the borehole using the electrical well-logging tool.\nThere is further provided a method for generating a log of fine features of a borehole wall with an electrical well-logging survey tool that is suspended from a cable inside the borehole which penetrates an earth formation, the method comprising the steps of:\nproviding a relatively non-conductive fluid, wherein the non-conductive fluid is selected from the group consisting of a non-aqueous fluid and an invert emulsion of an aqueous phase in a non-conductive liquid, wherein the aqueous portion of the invert emulsion ranges up to about 70 percent by volume, preferably from about 10 to about 40 percent by volume;\nadding an effective amount of a salt to the non-conductive fluid sufficient to modify the electrical characteristics thereof resulting in a modified fluid such that an electrical well-logging tool when disposed in the modified fluid can perform an electrical well-logging survey of a borehole containing the modified fluid,\nwherein the salt has a cation and an anion, the cation is selected from the group consisting of divalent and trivalent metal cations capable of forming a complex, and the anion is one that disassociates from the cation in an aqueous medium;\nproviding a column of the modified fluid in a portion of the borehole;\ndisposing an electrical well-logging survey tool in the portion of the borehole; and\nperforming an electrical well-logging survey of the portion of the borehole using the electrical well-logging tool, wherein the survey includes at least\ngenerating signals which represent a high spatial resolution measurement of a characteristic of the borehole and in the aggregate represent comparable high resolution features of the characteristic over an effectively continuous vertical and circumferential segment of the borehole wall and\ngenerating high resolution depth signals representative of the of the borehole depth to which the high spatial resolution characteristic signals relate;\nconverting the characteristic signals with the depth signals to produce the characteristic signals as a function of borehole depth;\nderiving from the latter signals, signals which represent the characteristic as a linear function of borehole depth;\ngenerating color, e.g., grey, scale values of the latter signals for display of fine features of the characteristic; and\nforming from the color scale values a visual image of the effectively continuous segment of the borehole as a linear function of borehole depth wherein the image has a color scale with which fine features of the characteristic of the segment of the borehole wall are visually enhanced.\nThere is also provided a method for generating a log of fine conductivity features of a borehole wall with a tool that is suspended from a cable inside a borehole which penetrates unearth formation and which tool has a plurality of survey electrodes arranged in multiple rows, preferably two rows, in a predetermined overlapping manner to investigate a circumferentially overlapping continuous segment of the borehole wall with high spatial resolution, the method comprising the steps of:\nproviding a relatively non-conductive fluid, wherein the non-conductive fluid is selected from the group consisting of a non-aqueous fluid and an invert emulsion of an aqueous phase in a non-conductive liquid, wherein the aqueous portion of the invert emulsion ranges up to about 70 percent by volume, preferably from about 10 to about 40 percent by volume;\nadding an effective amount of a salt to the non-conductive fluid sufficient to modify the electrical characteristics thereof resulting in a modified fluid such that an electrical well-logging tool when disposed in the modified fluid can perform an electrical well-logging survey of a borehole containing the modified fluid,\nwherein the salt has a cation and an anion, the cation is selected from the group consisting of divalent and trivalent metal cations capable of forming a complex, and the anion is one that disassociates from the cation in an aqueous medium;\nproviding a column of the modified fluid in a portion of the borehole;\ndisposing an electrical well-logging survey tool in the portion of the borehole; and\nperforming an electrical well-logging survey of the portion of the borehole using the electrical well-logging tool, wherein the survey includes at least generating high spatial resolution conductivity signals whose amplitudes represent the conductivity of the borehole wall opposite the survey electrodes with the conductivity signals in the aggregate representing conductivity of an effectively continuous vertical and circumferential segment of the borehole wall;\nconverting the conductivity signals to high spatial resolution conductivity signals which represent the conductivity of the borehole wall segment as a linear function of borehole depth; and\nforming from the latter conductivity signals a visual image of the effectively continuous borehole segment as a linear function of borehole depth and with a color, e.g., grey, scale with which fine conductivity features indicative of stratigraphy, vugs and fractures present in the borehole wall segments are visually enhanced."}
{"text": "1. Field of the Invention\nThe present invention relates to a heat exchanger which can be used, for example, in the freezing chamber of a refrigerator and, more particularly, to a heat exchanger having an enhanced cooling power to enable the freezing chamber to serve as a rapid freezing chamber.\nIn a refrigerator of the type broadly used, e.g., the chilled air circulation type, a cooling chamber is disposed at the innermost portion of a freezing chamber. The cooling chamber accommodates various parts such as a primary cooler, blower and so forth. In operation, air is forcibly supplied by the blower to the primary cooler and the air chilled by the primary cooler is discharged from the chilled air outlet provided at the upper part of the cooling chamber. The chilled air is introduced through a chilled air inlet provided at the lower part of the freezing chamber into the cooling chamber and is blown again from the chilled air outlet. In this type of refrigerator, a considerably long time is required for the cooling because the freezing chamber is cooled solely by the circulation of the chilled air. In fact, it takes about two hours for the water of normal temperature in a freezing pan placed in a refrigeration chamber to be frozen into ice.\nVarious improvements have been proposed up to now for shortening the cooling time. For instance, Japanese Utility Model Publication No. 52-52684 proposes to provide a heat exchanger, which is referred to as a \"secondary cooler\", in addition to the primary cooler mentioned above. This secondary cooler consists of a hollow vessel composed of a horizontal portion and a riser portion and filled with a refrigerant. The riser portion is held in contact with the primary cooler while the horizontal portion constitutes the bottom of a rapid freezing chamber. This type of refrigerator, however, can provide only a small heat exchanging efficiency because the area of contact between the refrigerant and the wall of the hollow vessel is limited. In addition, it is difficult to absorb the heat from the cooling object and to discharge the absorbed heat at a high efficiency because the horizontal portion for absorbing the heat and the riser portion for discharging the heat have an identical structure."}
{"text": "Given the complexity of computer systems, data, and the processing of data, challenges may arise for a computer storage system in attempting to differentiate between specific data types. In addition, computer storage systems are often unable to differentiate between classes of data or metadata due to limitations on the type of interface (e.g., a block interface), and the lack of a specific data-distinguishing mechanism. Thus, the computing performance may be negatively impacted.\nReference will now be made to the exemplary embodiments illustrated, and specific language will be used herein to describe the same. It will nevertheless be understood that no limitation on disclosure scope is thereby intended."}
{"text": "Third and fourth generation mobile telecommunication systems, such as those based on the 3GPP defined UMTS and Long Term Evolution (LTE) architecture are able to support more sophisticated services than simple voice and messaging services offered by previous generations of mobile telecommunication systems. For example, with the improved radio interface and enhanced data rates provided by LTE systems, a user is able to enjoy high data rate applications such as mobile video streaming and mobile video conferencing that would previously only have been available via a fixed line data connection. The demand to deploy third and fourth generation networks is therefore strong and the coverage area of these networks, i.e. geographic locations where access to the networks is possible, is expected to increase rapidly.\nThe anticipated widespread deployment of third and fourth generation networks has led to the parallel development of a class of devices and applications which, rather than taking advantage of the high data rates available, instead take advantage of the robust radio interface and increasing ubiquity of the coverage area. Examples include so-called machine type communication (MTC) applications, which are typified by semi-autonomous or autonomous wireless communication devices (i.e. MTC devices) communicating small amounts of data on a relatively infrequent basis. Examples include so-called smart meters which, for example, are located in a customer's house and periodically transmit information back to a central MTC server data relating to the customers consumption of a utility such as gas, water, electricity and so on. Other examples include medical devices which are continuously or intermittently transmitting data such as for example measurements or readings from monitors via a communications network to a server, and automotive applications in which measurement data is gathered from sensors on a vehicle and transmitted via a mobile communications network to a server attached to the network.\nWhilst it can be convenient for a terminal such as an MTC type terminal to take advantage of the wide coverage area provided by a third or fourth generation mobile telecommunication network, there are at present disadvantages and challenges to successful deployment. Unlike a conventional third or fourth generation communications device such as a smartphone, an MTC-type terminal is preferably relatively simple and inexpensive, having a reduced capability. In addition MTC-devices are often deployed in situations that do not afford easy access for direct maintenance or replacement, so that reliable and efficient operation can be crucial. Furthermore, while the type of functions performed by the MTC-type terminal (e.g. collecting and reporting back data) do not require particularly complex processing to perform, third and fourth generation mobile telecommunication networks typically employ advanced data modulation techniques (such as 16QAM or 64QAM) on the radio interface which can require more complex and expensive radio transceivers to implement.\nIt is usually justified to include such complex transceivers in a smartphone as a smartphone will typically require a powerful processor to perform typical smartphone type functions. However, as indicated above, there is now a desire to use relatively inexpensive and less complex devices to communicate using LTE type networks. Accordingly such devices may be low power or battery operated and therefore have a reduced transmission power capability compared with more conventional devices. Such MTC devices may also be deployed in remote locations where radio propagation conditions may be poor so that radio signals transmitted by the MTC device may be less likely to be receivable by a base station."}
{"text": "The following background information may present examples of specific aspects of the prior art (e.g., without limitation, approaches, facts, or common wisdom) that, while expected to be helpful to further educate the reader as to additional aspects of the prior art, is not to be construed as limiting the present invention, or any embodiments thereof, to anything stated or implied therein or inferred thereupon.\nGenerally, a skateboard is a type of sports equipment or toy used primarily for the activity of skateboarding, It usually consists of a specially designed Maplewood board combined with a polyurethane coating used for making smoother slides and stronger durability.\nTypically, skateboards are equipped with steering mechanisms known as trucks. The trucks are mounted on the underside of a skateboard opposite to each other, one in the front and one, in the rear. Each truck carries two wheels, one at. each end of the truck's axle. Most skateboards are sold as separate elements, namely the deck, trucks, and wheels. These elements are assembled. together with a few accessories, by either the buyer or the retail seller.\nPractically all trucks for skateboards, whether for serious competitive sport or pleasure use, comprise a mounting plate, known as a base plate, for bolting to the underside of the deck, a socket which may contain a bearing bush or liner and is formed in a block known as a yoke molded integrally with the base plate and depending therefrom as an angled projection, and a little spindle arm known as a hanger which is provided with a fixed axle for rotatably mounting one of the wheels at each end, The wheels of the skateboard are suspended in pairs from the underside of the deck by means of the trucks.\nTypically in operation, the skateboard is moved by pushing with one foot while the other remains on the board, or by pumping one's legs in structures such as a bowl or half pipe. A skateboard can also be used by simply standing on the deck while on a downward slope and allowing gravity to propel the board and rider.\nOther proposals have involved skateboards. The problem with these skateboards is that they do not magnify the force applied around the hose, and they are not self-powered. Also, the deck does not create a wave-like riding sensation. Even though the above cited skateboards meets some of the needs of the market, a bifurcated spring-tethered electric skateboard. The skateboard is bifurcated into a front deck and rear deck with each deck having independently powered trucks and wheels; whereby the skateboard has a mounted power source, generating electrical power for a motor that is transmitted to the wheels; whereby the skateboard also provides a controller that controls the speed of the motor, and thereby the wheels, is still desired."}
{"text": "In LED tubular lamps (lamps with a tubular outer contour which typically have several light emitting diodes (LEDs) as their light sources) space is very limited. As a result of the necessary electronics (drivers) for controlling the LEDs, in addition the lateral distance between a rest for applying an LED support populated with the LEDs and the LEDs is critical. In addition, other optics for light distribution, in particular light dispersion, are required in order to avoid inhomogeneous light distribution, which is caused by using the LEDs as point light sources."}
{"text": "Tray attachments for chairs, especially garden and patio chairs, are known. These attachments, however, require the original arm or arms to be modified by the provision of keyhole-type slots (U.S. Pat. No. 5,038,451 to Smith), which is beyond the skill and equipment of an ordinary person and makes the attachment non-retrofittable, or (U.S. Pat. No. 4,798,413 to Capelli) by drilling holes into the armrest or providing a specially modified armrest. U.S. Pat. No. 5,649,737 to Behnke proposes either the use of cloth-peg type fasteners or, again, the provision of holes or keyhole-type slots in the chair arms. U.S. Pat. No. 5,344,115 to Mayne, while not requiring modification of the original armrest, proposes the use of quite intricate, heavy, spring-loaded and thumbscrew-operated clamps that are not only expensive, but also certain to mar the polished surfaces of plastic patio or garden chairs."}
{"text": "This invention relates to treatment of catalyst supports.\nSupported chromium oxide catalysts have become an important vehicle for polymerization of 1-olefins such as ethylene and predominantly ethylene copolymers. As originally commercialized, these polymerizations were carried out under solution conditions. However, it was early recognized that a more economical route to such polymers was to utilize a slurry system sometimes called a particle form system, wherein the polymerization is carried out at a low enough temperature that the resulting polymer is insoluble in the diluent. A great deal of technology has been developed on means to adapt catalysts to compensate for the inherent difficulty in slurry polymerization systems of achieving high enough flow in the resulting polymer. However, another polymer property which is important in many applications is the molecular weight distribution which has an effect on the response to shear. There are a substantial number of applications which require a higher shear response polymer than is normally produced with chromium catalysts in a slurry polymerization system."}
{"text": "The use of control units such as shooting pots to introduce thermoplastic resin into single cavity molds is well known. In such arrangements a primary resin source feeds the reservoir of a shooting pot and the shooting pot, in turn, is operated to feed a measured or metered quantity of compound into a single cavity mold to form large items such as buckets, trash containers and like receptacles.\nIn contrast, the present invention is directed to the concept of supplying thermoplastic molding compound to a multi-cavity mold there the cavities of the mold are fed by a hot runner system where the hot runner system includes hooting pots (metering means or control units) which are individual to each cavity. That is, conduits that lead from the hot runner system to feed the respective mold cavities each contain a shooting pot or metering means for controlling precisely the introduction of a single shot of compound to a mating cavity on an individual or cavity by cavity basis.\nThe hot runner system is supplied or fed with thermoplastic material from a primary source usually in the form of an extruder or similar device in well known fashion.\nFor purposes of claiming the present invention, the shooting pots are referred to as metering means or as secondary compound advancing units.\nIt is a primary feature of the present invention to provide a hot runner system cooperating with a multi-cavity mold where the hot runner system includes metering means individual to each mold cavity.\nIt is a further feature of the invention to provide metering means in the hot runner system which are operable to feed precisely measured quantities of compound to each cavity.\nIt is a further feature of the invention to provide metering means individual to each cavity where the metering means are adjustable to change volume, for example, on a cavity by cavity basis. That is, each metering means may be adjusted individually without disabling the setting, function and operation of the other metering means in the hot runner system.\nA further feature of the invention is the treatment of a multi-cavity mold and a hot runner system as a unitary assembly.\nA still further feature of the invention is the provision of a novel process for feeding, metering and adjusting the introduction of thermoplastic compound into individual cavities of a multi-cavity mold by incorporating the feeding, metering and adjusting steps into a hot runner system.\nAn arrangement embracing certain features of the present invention may comprise, in combination, a multi-cavity mold and a hot runner system, supply means incorporated in the hot runner system for feeding measured quantities of thermoplastic compound to each said mold cavity on an individual cavity basis.\nA method embracing certain other features of the present invention may comprise the steps of providing a mold having a plurality of cavities, providing a hot runner system for feeding said cavities, connecting the hot runner system to each cavity by means of compound conduits individual to each cavity and placing metering means in each said cavity conduit for controlling the introduction of compound to each cavity on an individual cavity basis.\nOther feature and advantages of the present invention will become more apparent from an examination of the succeeding specification when read in conjunction with the a pended drawings in which;"}
{"text": "1. Field of the Invention\nThe present invention relates to applications of modular computer cores. More particularly, the present invention relates to the use of a hybrid circuit including a high-power processor and a low-power processor to provide a selectable “thermal design power” (TDP) in a thermal module together with an air-flow design.\n2. Discussion of the Related Art\nIn recent years, small and thin computing devices are highly favored. Some examples of small and thin computing devices include the iPad and the iPhone from Apple Computer, Inc., the “ultrabooks” notebook computers from Intel Corporation and its partners, and the ultra-thin “Android smartphones” from Google, Inc. and its partners. To support these “ultra” devices, microprocessor manufacturers have provided low-power microprocessors (e.g., the ARM microprocessors, or fan-less x86 microprocessors). These microprocessors—which dissipate less than 3 W TDP (i.e., 3 watts of “Thermal Design Power”)—are primarily targeted for basic applications. To execute more advanced applications, higher power microprocessors are needed. But microprocessors that have a TDP that is higher than 3 watts require a proper thermal module for heat dissipation. In addition, the size of a typical ultra-computing device (i.e., roughly, the size of a smartphone) makes it difficult to squeeze a proper thermal module into the limited space. Other constraints on such a device include: (i) the device as a whole has to be as light as possible; (ii) the form factor has to be handheld size; (iii) noise has to be kept to a minimum, so as to be non-intrusive on the user carrying it very closely to the body (e.g., in a shirt pocket); and (iv) the exterior case temperature has to be kept low enough to be handheld permissible. Therefore, a new thermal module that is able to dissipate heat of a high-power microprocessor within the limited space of an ultra-computing device is desired to support advanced applications.\nMetal blocks have been used as heat sinks that are mounted on low-power fan-less microprocessors to dissipate heat. It is also common to use metal chassis or cases to serve as passive heat sinks for low-power microprocessors. However, to transfer heat away from a higher power microprocessor, a much larger and more complete thermal module is required. Such a thermal module may include a heat dissipation plate, a heat pipe, and a heat sink. Further, it is customary also to include an integrated fan to increase airflow over the heat sink to expel the heated air out of the chassis or case quickly.\nExcluding the display and the touch panel, the body of a typical ultra-thin device is less than one centimeter thick. A bulky heat dissipation block does not fit in this thickness. In addition, it is impossible to put a conventional cooling fan within the confines of the smartphone-size computer. It is a challenge to computer supplies to find a design that cools down a smartphone size computer in which a high-power microprocessor is used. In an attempt to provide such a solution, some computer thermal management companies (e.g., SUNON in Taiwan) designed powerful “mighty mini-fans” that fit into the limited space. However, these new “mini” products do not generate enough airflow to cool a high-power microprocessor in an effective manner."}
{"text": "PWM-input, class-D amplifiers that do not require an external clock are known as conventional amplifiers (refer to Patent document 1, for example).\nIn the amplifier disclosed in Patent document 1, a difference signal between a negative feedback signal from a power stage and a PWM input signal passes through two integrators and is thereby demodulated sufficiently. PWM modulation is realized without using an external clock by comparing, with a comparator, the sufficiently demodulated difference signal with a clock signal that is taken out of the first integrator and in which PWM harmonics remain."}
{"text": "1. Field of the Invention\nThe present invention relates an erroneous operation protection circuit for preventing an erroneous operation due to coupling noises in a large scale integrated circuit (LSI) in which data bus wirings are aligned in close vicinity to each other on a semiconductor chip, and a large scale integrated circuit (LSI) into which the erroneous operation protection circuits are incorporated.\n2. Description of the Prior Art\nIn recent years, according to the progress of fine pattern technology and the request for improvement in system performance, the large scale integrated circuit (LSI) is manufactured on a larger scale and is operated at a higher speed. FIG. 1 shows an example of a microprocessor in which a plurality of circuit blocks 101, 191 are aligned on a semiconductor chip (LSI chip) 100. As shown in FIG. 1, a data bus wiring 110 to which the circuit blocks 101, 191 are connected is disposed between a plurality of sender (transmitter) side circuit blocks 101 and a receiver side circuit block (or the next stage circuit block) 191. This kind of the data bus wiring 110 is provided to correspond to a bit width which is processed simultaneously by the LSI chip 100 and provided in parallel at as narrow pitches as possible to reduce a wiring area. A wiring length of the data bus wiring 110 becomes long since the data bus wiring 110 extends over a plurality of circuit blocks 101, 191 on LSI chip.\nConversely, in order to improve the device packing density on a single LSI chip, wiring widths and wiring intervals of the metal wirings in the LSI chip are reduced year after year. However, a certain lower limit value is set to thicknesses of the metal wirings because the wiring resistance must be maintained small. In other words, the thicknesses of the metal wirings cannot be reduced below the lower limit value and therefore only an interval between both metal wirings is reduced while mutually opposing areas are kept at a certain value. Thus, there is a tendency to increase capacitance between mutually neighboring data bus wirings. Therefore, a rate of capacitance (wiring capacitance) between adjacent data bus wirings to a total wiring capacitance of the LSI is increased. In this manner, the capacitance between the two adjacent data bus wirings in the LSI corresponding to adjacent bits is extremely high and therefore influence of the coupling (coupling noises) between data on the concerned adjacent bits becomes large when the potential level of one of the adjacent bits is swinged.\nIn the prior art, as a system driving the data bus wirings, there has been known \"a CMOS system\" in which the potential level of respective wirings are swinged between a power supply level (\"H\" level) and a ground level (\"L\" level) with the use of circuits made of n-MOSFETs and p-MOSFETs being connected in series between the power supply and the ground. As another system driving the data bus wirings, \"a precharge type system\" in which respective wirings are set to the power supply level beforehand and then the n-MOSFETs arranged in the circuit block are turned ON only when the signal propagates the data bus at the ground level is also known. The precharge type system has such excellent features rather than the CMOS system that (1) no Miller effect is caused, and (2) direction of data transition is constant and thus optimization for the transition direction can be implemented in circuit design.\nFIG. 2 is a circuit diagram showing an example of the precharge type system, or a configuration of precharge type data bus wirings, which are generally used as the data bus wirings in the prior art. In FIG. 2, data bus wirings 110n-1, 110n, 110n+1 to which a plurality of circuit blocks 101 are connected have a function of propagating output data from the circuit blocks 101 to input buffers 111 of the circuit block 191 respectively. In respective circuit blocks 101, n-MOSFETs 102 for driving the data bus wirings 110n-1, 110n, 110n+1 into an \"L\" level are provided. Precharge circuits 112 for setting respective data bus wirings into an \"H\" level beforehand and latch circuits 113 for holding the data when the data bus wiring is at the \"H\" level are connected to the data bus wirings 110n-1, 110n, 110n+1.\nIn turn, an operation of the circuit shown in FIG. 2 will be explained. The data bus wirings 110n-1, 110n, 110n+1 are set to the \"H\" level by the precharge circuits 112 in advance. When one of the plurality of circuit blocks 101 transfer the data via the data bus wirings 110n-1, 110n, 110n+1, respective precharge circuits 112 are turned OFF. When the circuit block outputs the \"L\" level, the circuit block 101 for transferring the data renders the n-MOSFET 102 to turn ON and thus drive the corresponding data bus wiring into the \"L\" level. On the contrary, when one of the circuit blocks 101 outputs data at the \"H\" level, the latch circuit 113 continues to hold the \"H\" level of the corresponding data bus wiring while the corresponding n-MOSFET 102 in the circuit blocks 101 is kept in its OFF state. Thereby, the concerned data bus wiring can propagate the \"H\" level. As a consequence, the input buffers 111 of the circuit block 191 positioned on the preceding stage can receive desired data OUT.sub.n-1, OUT.sub.n, OUT.sub.n+1.\nHowever, if the integration density of the LSI is improved, as described above, capacitances 114 between the two adjacent wirings are increased in the circuit shown in FIG. 2 and thus capacitances between the two adjacent wirings are increased rather than the capacitance between the data bus wirings 110n-1, 110n, 110n+1 and the ground level. In this case, even if the LSI circuit is designed such that one of the data bus wirings is driven into the \"L\" level and the adjoining data bus wiring is held at the \"H\" level by the latch circuit 113, the adjoining data bus wiring to be held at the \"H\" level is shifted into the \"L\" level due to influence of the coupling therebetween. As a result, data at the \"L\" level is transferred to the gate in the circuit block 191 to cause the erroneous operation.\nAs the countermeasure for such disadvantage, there is a method wherein transition to the \"L\" level can be prevented by increasing a driving force of the latch circuit 113 rather than the force generated by the influence of the coupling between two adjacent data bus wirings. However, in this case, when the circuit block 101 outputs the \"L\" level onto the data bus wirings, the latch circuits 113 drive the data bus wirings strongly into the \"H\" level. As a result, there is such a disadvantage that collision of data happens and the delay of data is caused.\nAs a method of preventing erroneous operation by avoiding such collision of data, as shown in FIG. 3, a method has been proposed wherein the data bus wirings 110n-1, 10n, 10n+1 are driven by tri-state buffers 103 in respective circuit blocks 101. The tri-state buffers 103 can drive respective data bus wirings 110n-1, 10n, 10n+1 into the \"L\" level, the \"H\" level, or its floating state. In the example shown in FIG. 3, when the data bus wirings 110n-1, 10n, 10n+1 are precharged or when the circuit blocks are inactivated, the tri-state buffers 103 are brought into their non-output state not to drive the data bus wirings 110n-1, 10n, 10n+1. Only when the circuit blocks are activated to output the data, the tri-state buffers 103 can drive the data bus wirings into the \"L\" level or the \"H\" level according to the data.\nAs in the above circuit example shown in FIG. 2, respective data bus wirings 110n-1, 10n, 110n+1 are set to the \"H\" level in advance and data on the data bus wirings are shifted only when the data is at the \"L\" level, nevertheless the data bus wirings to be held at the \"H\" level are driven into the \"H\" level by the tri-state buffers 103 in the circuit blocks 101 which are activated. Hence, even if data transition is caused due to the coupling between the adjacent data bus wirings, the tri-state buffers 103 drive strongly the data bus wiring to the \"H\" level. Therefore, transition of data to the \"L\" level can be suppressed and also erroneous operation can be prevented.\nHowever, in the circuit shown in FIG. 3, the tri-state buffers, each having a driving force larger than the force generated by the influence of the coupling, must be provided according to the number of the circuit blocks. This is because such tri-state buffers must be arranged in or near respective circuit blocks since the tri-state buffers need the output signals of respective circuit blocks. In the event that the tri-state buffers are connected to the data bus wirings according to the number of the circuit blocks, the number of transistors constituting the tri-state buffers is of course increased. Therefore, junction capacitances of the transistors are increased as a whole and total parasitic capacitances of the data bus wirings are increased. As a consequence, stray capacitance of the data bus wirings are increased so that wiring delay is increased. In addition, in case the length of the data bus wirings are extremely prolonged to the extent that special variation in the signal level is caused due to the wiring delay on the same data bus wiring, it is difficult to suppress the coupling noises at only one location on the data bus wiring.\nLike the above, in the prior art, there has been a problem that a large number of tri-state buffers must be prepared to suppress the coupling noises and thus stray capacitances of the data bus wirings are significantly increased to increase the wiring delay."}
{"text": "In a work vehicle, when an engine load is equal to or lower than a rated engine torque, the engine torque is matched to the engine load in a high-speed control area in a torque chart. For instance, a target engine speed is set according to the setting of a fuel dial and the high-speed control area associated with this target engine speed is set.\nAlternatively, the high-speed control area is set according to the setting of a fuel dial and the target engine speed associated with this high-speed control area is set. The engine load and the engine torque are matched in this high-speed control area.\nMany operators generally set a target engine speed at or around a rated engine speed so as to improve an operating quantity. A low fuel-consumption area, namely a fuel-efficient area, usually exists in a middle-speed area or a high-torque area on an engine torque chart. Therefore, a high-speed control area defined between a non-load high-idle speed and a rated speed does not correspond to an efficient area in terms of fuel consumption.\nIn order to drive an engine in the fuel-efficient area, a typically-known control device presets the value of a target engine speed and the value of a target engine output torque, which values correspond to each other, for each of plural selectable operation modes (see, for instance, Patent Document 1). With the use of such a control device, when an operator selects, for instance, a second operation mode, the engine speed can be set lower than that in a first operation mode, and therefore the fuel consumption can be improved.\nHowever, according to the above-described operation mode switching, the operator needs to operate the operation mode switching each time to improve the fuel consumption. Further, in a situation where the engine speed in the second operation mode is set at a value simply reduced relative to the engine speed in the first operation mode, selection of the second operation mode leads to the following problem. The maximum speed of a working device (hereinafter referred to as a work equipment) of a work vehicle is decreased as compared to that in the first operation mode. As a result, an operating quantity in the second operation mode becomes smaller than that in the first operation mode.\n[Patent Document 1] JP-A-10-273919"}
{"text": "The present invention relates to an optical scanner, and particularly to the optical scanner by which the light is 2-dimensionally scanned.\nPresently, the optical scanner by which the ray such as the laser light is deflected and scanned, is used for an optical device such as a bar-code reader, laser printer, display. Conventionally, as the optical scanner such as this, there is a polygonal mirror in which a polygonal column mirror is rotated by a motor, and a galvano-mirror in which a plane mirror is rotation-oscillated by an electromagnetic actuator. However, in a mechanical structure which is driven by the motor or electromagnetic actuator, because the shape of the structural parts is large or cost is high, there is a problem that the downsizing of the device for which the optical scanner is used, is hindered, or results in high cost. Further, when the ray is 2-dimensionally scanned, generally, a device in which the polygonal mirror and the galvano-mirror are combined is used, however, in order to conduct the accurate 2-dimensional scanning, it is necessary that the respective scanning directions are accurately positioned so that the scanning directions by respective mirrors are orthogonal with each other, and there is also a problem that the optical adjustment is very complicated.\nTherefore, in order to realize the downsizing, low cost, and the increase of the productivity, by using a micro-machining technology by which silicon or glass for which the semiconductor manufacturing technology is applied, is micro-machined, various micro optical scanners in which the structural parts such as the mirror or an elastic beam are integrally formed on the semiconductor substrate are successively developed.\nFor example, two frames of an inner frame for supporting the mirror through the first double holding beam, and an outer frame for supporting the inner frame through the second double holding beam are provided, the first and the second double holding beams are orthogonal with each other, and the bending-ness of the second double holding beam and the resonance frequency of the oscillation of the first double holding beam, are made close. Then, the magnetic distortion film is formed on the second double holding beam, and the engineering of the optical scanner (refer to Japanese Patent Unexamined Publication: Tokkai No. 2003-207737) by which, when the distortion and the bending oscillation are generated, the mirror is rotated by using the first and second double holding beams as 2 axes, and 2-dimensionally optical scanning is conducted, further, in almost the same structure] as the optical scanner disclosed in Patent Document 1, 2 frames and 2 double holding beams are provided, and by a plurality of piezoelectric elements arranged in the inner frame and the outer frame, when the rotation torque is respectively independently, acted on the inner frame and the mirror, the engineering of the optical scanner (refer to Tokkai-No. 2005-148459) which rotate the mirror when the first and the second double holding beams are made 2 axes, is disclosed.\nHowever, the optical scanner disclosed in Patent document 1, the distortion and bending oscillation are simultaneously conducted on the second double holding beam by the actuator.\nThat is, because distortion and bending can not be independently conducted, it is considered to be difficult that the horizontal and vertical deflection are independently controlled. Further, because the resonance frequencies of the distortion oscillation of the first double holding beam and the bending oscillation of the second double holding beam are set so that they are almost same value, when the horizontal and the vertical deflection frequencies are largely different, that is, when the resonance frequencies of the distortion and the bending oscillation are largely different, it is considered to be difficult that the first double holding beam and the second double holding beam are resonated at respectively necessary deflection frequency. Further, the optical scanner disclosed in Patent Document 2, because, as for the horizontal deflection and the vertical deflection, a plurality of piezoelectric element are respectively independently provided, there is a possibility that it results in the complexity of the device, and high cost."}
{"text": "Field of the Invention\nThe invention relates to a method and system for producing screen or raster data for imaging units of a printing machine.\nPrinting machines with imaging units have become known heretofore. Imaging units of this type produce a printing image on a printing plate of the printing machine by scanning with a laser beam, the position and size of the screen or raster dots of the printing image being controlled by screen or raster data, with which the laser beam is modulated in the course of a scanning movement over the surface of the printing plate.\nWith the aid of imaging units of this type, as compared with the conventional film exposure technique, a plurality of processing stages on the path from an original image to the finished print are dispensed with, so that alternating print jobs can be processed speedily and cost-effectively. Nevertheless, even when such imaging units are used, the processing process from the original image to the finished print breaks down into a multiplicity of processing stages which have to be performed successively and which, to some extent, are time-consuming. In a conventional screen production system, as illustrated in FIG. 1 of the drawings herein, raw image data initially passes through a screen processor 1 (screen image processor RIP), which converts this data into a plurality of partial images, respectively, one partial image for each color to be printed, while simultaneously performing all of the necessary calibration operations. These partial images, respectively, are transferred by an administration unit 2 into a buffer storage unit 3. After the screen processor 1 has completely processed the original image, and all the necessary partial images are present in the buffer storage unit 3, they can be transferred therefrom to the imaging units 4 of a diagrammatically illustrated printing machine 5. Buffering the screen data in the unit 3 is necessary, because the imaging units 4, respectively, must be supplied with data without interruption at a data rate which is predefined or prescribed by the type of construction thereof, and a non-buffered imaging process will fail if the screen processor 1 is not capable of supplying the data at the required rate.\nConsequently, the imaging units 4 remain inactive while the screen processor 1 is operating, and the reverse. Because the screen data for each screen dot to be produced by the imaging units must specify the size of the respective dot, they include a significantly greater quantity of data than the raw data originally input into the screen processor 1, which, for example, can be a file in Postscript or PDF format. Storing the screen data for the situation wherein they are used again for a subsequent print job is therefore rather complicated. In addition, if the screen processor 1 performs a high-quality calibration, the dot increases or growths of the individual imaging units of the various stages of the printing machine are taken into account in the specification of the screen dots. These dot increases or growths can vary from one imaging unit 4 to another, within certain limits, because of the scatter in the optical properties of the laser and the scanning system, which are used. Consequently, the screen processor 1 produces the screen data for a given printing ink, respectively, taking specifically into account the properties of the imaging unit 4 which sets the image on the printing plate for the relevant color. If a print job is to be repeated based upon stored screen data, it is consequently necessary for the distribution of the printing inks to the various stages of the printing machine 5, at the time the print job is repeated, to be the same as that for which the screen data were originally calculated. If, in the interim, jobs with a different sequence of colors have been processed, this would require lengthy washing and re-inking of the printing machine.\nFor a repetition of the print job on a different machine, renewed screening calibrated to the imaging units of the machine is always necessary, even if the machines are of identical construction and are equivalent in all the other relevant parameters, such as the printing material, the color used, the plate material, and so forth.\nIt is accordingly an object of the invention to provide a method and system for producing screen data for imaging units of a printing machine, which permit further acceleration of operating steps from the original raw image to the print.\nWith the foregoing and other objects in view, there is provided, in accordance with one aspect of the invention, a method of producing screen data for imaging units of a printing machine, which comprises the steps of breaking down raw image data into a plurality of partial images, respectively, corresponding to one printing ink; outputting the partial images to a plurality of screen processors, respectively, corresponding to the number of colors to the printed; and processing the partial images simultaneously for forming screen data by the screen processors for output to a respective one of the imaging units.\nIn accordance with another mode, the method of the invention includes performing by the screen processors, respectively, a calibration of the partial images relative to dot growth of the imaging unit with which it is associated.\nIn accordance with a further mode, the method of the invention includes breaking down the raw image data in a preprocessing unit that is separate from the screen processors.\nIn accordance with an added mode, the method of the invention includes performing a calibration by the preprocessing unit, relative to properties of at least one of the printing material, the printing ink, the blanket cylinder, and the plate material of the printing machine.\nIn accordance with an additional mode, the method of the invention includes, by the preprocessing unit, outputting partial image data to the screen processors before the step of breaking down the raw image data has been finished.\nIn accordance with yet another mode, the method of the invention includes buffering the screen data before outputting the screen data to the respective imaging units.\nIn accordance with yet a further mode of the method invention, the step of breaking down raw image data includes a trapping procedure.\nIn accordance with a concomitant aspect of the invention, there is provided a screen production system for a printing machine having imaging units, comprising a preprocessing unit for receiving raw image data and for breaking down the raw image data into a plurality of partial images, respectively, corresponding to one printing ink, and a plurality of screen processors, respectively, for screening the partial images, respectively, and for outputting to the imaging units, respectively, partial-image screen data obtained from the respective partial images.\nThus, instead of determining, in order, the sizes of the screen dots of the colors to be printed for each point or dot in an image to be printed here, initially, a separation into various colors is performed, and the partial image data, respectively, corresponding to one color, is processed by a screen processor which is specifically provided for the purpose and which is therefore capable of providing the screen data of the associated color significantly more rapidly than heretofore.\nThe calibration of the partial images relative to the dot growth of the imaging unit, the calibration being necessary for a high-quality print, is preferably performed by the screen processor which is associated with the imaging unit. Those partial images not yet screened, which form data files of a relatively small size, can comfortably be stored intermediately for a repetition of the print job. When the print job is to be repeated, the intermediately stored partial images are screened anew. If, in the interim, the distribution of the colors to the stages of the printing machine has changed, or another printing machine is to be used, the renewed screening can be performed by a screen processor that differs from the first one and can be calibrated to the dot growth of a different printing stage.\nThe step of breaking down the raw image data into partial images is preferably performed in a preprocessing unit that is separate from the screen processors. This preprocessing unit can perform a calibration relative to all those printing parameters which are standard for the various stages of the printing machine, thus, for example, the printing material used, the printing ink, the blanket cylinder or the plate material of the printing machine.\nA particular advantage in using the preprocessing unit is the possibility of intermeshing in time the actions of breaking down and screening, by the preprocessing unit outputting partial image data to the screen processors before the step of breaking down the raw image data has been finished completely.\nIn order to be able to ensure a uniform data stream at the rate required by the imaging unit, the screen data are preferably buffered before being output to the imaging unit. However, such a buffer can be considerably smaller than that needed to store an entire partial image. If the computing power of the screen processors is sufficiently high so that they can ensure a stream of screen data at the rate required by the imaging unit, the buffer can even be omitted completely.\nThe step of breaking down the raw image data preferably also includes a trapping procedure.\nThe object of the invention is also achieved by a screen production system for a printing machine which includes a preprocessing unit for receiving raw image data and for breaking down the raw image data into a plurality of partial images, each of which corresponds to one printing ink, and a plurality of screen processors for screening one of the partial images, respectively, and for outputting the thus obtained partial-image screen data to an imaging unit.\nOther features which are considered as characteristic for the invention are set forth in the appended claims.\nAlthough the invention is illustrated and described herein as a method and embodied as a system for producing screen data, it is nevertheless not intended to be limited to the details shown, since various modifications and structural changes may be made therein without departing from the spirit of the invention and within the scope and range of equivalents of the claims.\nThe construction and method of operation of the invention, however, together with additional objects and advantages thereof will be best understood from the following description of specific embodiments when read in connection with the accompanying drawings, wherein:"}
{"text": "With the growing popularity of personal computers, a variety of applications address the needs of the graphics industry. Documents of all kinds are being converted to digital data which is manageable by electronic devices. Graphics has inundated the PC (personal computer) marketplace in the last ten years, due to the implementation of low cost, high-density memories, VLSI graphics controllers, graphics monitors, graphics standards and graphics applications software.\nThe PC user operates in a monochrome or color graphics platform, depending on the application. Cost used to be significant, and forced the low-end graphics user to select a monochrome system, but recent low-cost color systems have now virtually eliminated this limitation. Within the color marketplace, there are the standard color formats applicable to both the low-end mass marketplace and the high-end specialty marketplace.\nThe first category typically involves about $500 to $800 for a color graphics card and monitor, and the vast majority of color PC users fall into this category. The second category can involve between $6,000 and $10,000, and only specialty users can justify this cost. Currently, color scanning is available basically for the second of these categories, because it requires high spatial resolution (1024.times.768) and high-color resolution (24-bit true color) to handle the large data rates. It is a principle goal of the present invention to provide color scanning to the low-end color graphics user.\nRelevant popular applications of graphics include illustrations, presentations and publishing. These three applications were each based in large part on the availability of input data and output hard copy technology. The introduction of black and white scanners brought a personal touch to these three applications which was previously impossible. Even more dramatic was the introduction of low-cost laser printers. With these, a user can input data from conventional sources and output it to hard copy, all with an affordable desktop system.\nThe richness of color was predominantly ignored, due to the unavailability of color images and the expensive printing costs. Now, however, color printers are becoming available at affordable costs, and consequently color scanners are becoming popular. The enormous body of color photographs had remained untouched, due to the costs. This is changing, and those same applications which could afford black and white scanners and laser printers will soon be able to afford color processing. These users will want and need access to color images, and in particular will want to access the images stored on color slides.\nCurrently, the three graphics applications mentioned above rely on clip art, but current clip art libraries pale in comparison to the images available on 35 mm color slides. An inexpensive slide scanner is essential for users who wish to access these color images. A graphics user can currently get several clip illustrations on a single floppy disk, due to the fact that the binary image format is relatively non-data intensive. A 320 by 200 pixel black and white image can be stored in 8000 bytes. Unfortunately, this is not the case for color images. Currently, a 24-bit color image of the same size would require 192,000 bytes. At current digital storage densities, access speeds and costs, it would not be practical for graphics users to utilize color images in digital formats, especially when one considers that a full page image of 3600 by 2400 pixels requires about 26 megabytes whereas the maximum storage currently available on a standard floppy disk is only 1.44 megabytes.\nIn order to utilize these color images, it will be necessary to store these color images in their native format, namely on 35 mm color slides. Each slide is small and very inexpensive to reproduce and store. A single standard 120 slide carrousel could store over 3 gigabytes of information. Furthermore, the slides retain their data with minimal care. An inexpensive slide scanner is therefore an absolutely necessary component in the graphics user's computer system.\nThere are several color scanners currently available on the market which are capable of scanning transparencies, but none of these scanners address the needs of the low-end market. By far the greatest majority of PC users operate with 8-bit color graphics cards and desire a scanner which is simple to use, compact, reliable and inexpensive. A user would not be likely to spend as much money on a slide scanner as was spent on the entire computer system. In addition, it is highly desirable for the scanner to actually mount within the PC, to reduce desktop clutter. Although flexible and powerful, current transparency scanners are expensive, and are designed with a physical size which cannot practically be made sufficiently compact to fit within a conventional PC. Albeit, these allow a wide range of sample size formats (for example from a 35 mm slide transparency to an 8.5 by 11 inch opaque), but at a cost which is unrealistic to the end user. An object of the present invention is to provide a compact and inexpensive system especially designed for a particular size application (for example 35 mm transparencies).\nObjects of the present invention include provision of a scanner system which can generate a digital image from a transparency, which is compact, reliable, flexible and inexpensive, and which facilitates computer-based illustration, artwork, presentation graphics, desktop publishing and image database management.\nA further object is the provision in the system of capability to directly digitize, store and display images without the need for any intermediate image storage, such as in solid-state or magnetic storage media.\nA further object is the provision of capability to display an image with minimal delay, and preferably to effect a full color scan digitizing 480.times.640 pixels in less than ten seconds from the moment a slide starts to move into the system.\nIt is another object to allow a user to vary the magnification and the location of the image origin (zoom and pan capability), the zoom factor not being limited to integer multiples of the pitch of a solid-state image sensor and the zoom not being implemented in software.\nA further object is the provision of a mechanical chassis with the size and shape of a conventional floppy disk drive so that the scanner can be directly installed in a conventional personal computer frame, wherein transparencies or negative strips are inserted into the system through an opening in its front panel, and power requirements of the scanner are met by the internal power (5 Volt, 12 Volt, 20 Watt maximum) supplied in popular computer systems for use by peripherals.\nA further object is to provide a system which is self-calibrating, and does not require extensive warm-up periods or set-up procedures.\nA further object is the provision of a program for the host computer which controls the scanner and allows an operator to interact with the scanning process, and which preferably establishes a windowed environment which allows the operator to open and close files, edit existing images, modify the layout of the window environment, control the color coding, and control the scanner."}
{"text": "(1) Field of the Invention\nThe present invention relates to the field of sonar sensors and in particular to towed, active sonar arrays.\n(2) Description of the Prior Art\nAlthough passive (receiving) towable sonar arrays have become refined technology, active (transmitting) arrays, because of the requirement to position the array in an upright, vertical position, have presented certain drawbacks. Typical conventional active systems are mounted in hard, typically large, rigid towed bodies, for example, the current AN/SQS-35 tow body. The tow body is required to maintain the active array in a proper vertical orientation in order to provide the required acoustic pattern. Conventional towed active sonar systems are large, cumbersome, and require expensive handling systems that use a substantial amount of space aboard a surface combatant. Depending on the frequency of the transducers, the hard, towed bodies are quite large and difficult to handle, frequently weighing up to 4000 pounds or more. Additionally, the handling equipment needed to deploy and recover such a tow body requires considerable deck space on the aft end of the ship and this equipment presents a large radar target. The launch doors, chutes, and associated hardware also increase the radar signature of the ship. These types of handling systems impose limitations on the maneuverability of the tow vessel, prevent the covert deployment of the active array, and are unwieldly and dangerous for deployments or recoveries during high sea states. A means is needed whereby an active transducer array can be deployed from a ship in a manner similar to the deployment of passive transducer arrays, such as paying out through a hull-mounted port. The active transducer array must, nevertheless, tow in a substantially vertical position after deployment."}
{"text": "This invention relates to a quick-disconnect coupling arrangement for connecting two members of a linkage between which respective longitudinal, lateral and rotational movement is to be prevented. The prior art recognizes that, where linkage members are attached to a structure, it may be inconvenient or impossible to disconnect the linkage from that structure when removal of that structure is necessary or otherwise. In such a case, it is preferable to use two linkage members rigidly but releasably connected so as to prevent lateral or rotational movement therebetween. The most common form of such a linkage is two bars with respective left and right-hand threads and a connecting turnbuckle. Turnbuckles enjoy the significant advantage that their installation, removal and adjustment requires no special tools. The turnbuckle has its disadvantages, however, in certain applications. In order to loosen and remove a turnbuckle coupling, the turnbuckle must be unscrewed, which forces the right and left-hand threaded end portions of the respective linkage members apart longitudinally. This operational characteristic renders turnbuckles inappropriate for use where longitudinal or rotational movement of the linkage members is not possible or where it is preferable to maintain longitudinal adjustment of the linkage members. An example is vehicles having a removable machine element such as a power package, transmission, or the like, with a linkage for attachment to another machine element, such as a shifter, throttle, clutch, or the like. Removal of the removable machine element requires its disconnection from the other machine element. Longitudinal and rotational movement of the linkage members and disconnection of the members from either machine element may be impossible. Disturbing the longitudinal adjustment of the linkage is likely to be inconvenient.\nUse of clevis and eye connections including a connecting pin are commonly known in the prior art. A standard pin, clevis and eye arrangement prevents longitudinal movement between adjacent linkage members, but allows relative rotation of such members about the pin. Use of a clevis pin has certain recognized disadvantages, however. The pin must be maintained in position by means of a cotter pin, retaining ring, or the like. Cotter pins tend to break, fall off, or get lost. Moreover, their application and removal requires a pair of pliers. Likewise, installation and removal of retaining rings requires a special tool. The prior art has overcome the disadvantages associated with cotter pins and retaining rings by using spring pin connectors. These connectors include a pin for insertion through the eye and clevis holes, a snap for attaching the connector to a neck on the pin or clevis, and a strap connecting the pin and snap."}
{"text": "Winches are devices used for traction, and are mainly used to rescue vehicles in harsh environments. When the vehicle is trapped in muddy or rough terrain and not the wheels are on the ground, the winch can be used to drag the vehicle out by using another object as a fulcrum. The winch can also be used to drag large-size obstacles such as logs. Existing electrical winches are mainly comprised of a motor, a driving shaft, a gear reduction box and a winding drum. The gear reduction box transmits the driving power from the motor to the winding drum so that a substantial traction force can be induced on a steel wire wound around the winding drum for functions of the winch.\nWinches are mainly prepared for emergencies so they have to work under various hostile environments. Unexpected failures may frequently happen during use of winches. The structures of existing winches need improvement in ease of maintenance and operation to deal with unexpected failures rapidly. Besides, mechanical components of a winch are subject to water damage in water or rainy weather."}
{"text": "The subject matter disclosed herein generally relates to storage of materials and reagents used in biological and chemical processes.\nBiologically active materials may be used in a variety of laboratory and analytic contexts. However, in general, such biologically active materials may have a relatively short shelf life if not treated or prepared to enhance their storage characteristics.\nFor example, various approaches for biological reagent stabilization and storage are known. One such approach includes storing the protein in a liquid format at reduced temperature (e.g., −20° C. to 8° C.). For instance, certain biological reagents may be stored in a 50% glycerol solution maintained at 4° C. or as low as −20° C. Alternatively, certain biological reagents may instead be frozen for storage, such as at or below −20° C. Obviously both of these storage approaches require refrigeration to maintain the biological activity of the reagent for extended time.\nIn addition, biological reagents may be stored in a lyophilized form, in which the reagent is dried by freezing in a high vacuum. Such lyophilized reagents may be stored at low temperature or at room temperature. However, the processes used to produce the lyophilized product may be complex and time consuming. In particular, certain such processes used to produce lyophilized cakes, films, beads or spheres of biological enzymatic mixtures that may be batch processes, that do not allow for continuous production of the product. In some methods, frozen solutions are dehydrated, requiring a complicated freeze drying method. Further, in such techniques the desired manufacturing tolerances, such as with respect to bead size, may be difficult to obtain or maintain."}
{"text": "Lasers, amplifiers and other optical apparatus based on optical fiber can provide flexible, rugged and relatively simple sources of optical energy. Accordingly, in many applications such optical fiber apparatus can often have one or more advantages as compared to counterparts based on a gas medium (e.g., CO2) or on a bulk solid-state medium (e.g., a Nd:YAG rod). For example, optical fiber lasers can have a smaller footprint, or can be more efficient, or can require less sophisticated cooling arrangements as compared to using a gas or bulk rod solid-state laser in a similar application. Often, however, it can be desirable to increase the output power of optical fiber apparatus, as certain gas and bulk solid-state lasers can readily produce high CW output powers or pulses of optical energy having high energy and/or high peak power.\nUnfortunately, because of the high power density inherent in confining optical energy to the relatively small cross sectional area of an optical fiber, non-linear phenomena, such as Stimulated Raman Scattering (SRS) or Stimulated Brillouin Scattering (SBS), can severely limit scaling the output power of a fiber laser or amplifier to higher powers. Though these non-linear processes are complex, each can be addressed, at least in part, by reducing the power density in the core of the fiber. One way to reduce power density is to increase the diameter of the core of the fiber and/or reduce the numerical aperture (NA) of the core, such that the fiber has a larger mode field diameter (MFD). Reducing the power density in this manner can increase the power threshold for the onset of the undesirable non-linear phenomena.\nThis approach, however, is not without drawbacks. Fibers having larger core diameters can typically support higher order transverse modes (e.g., LP11, LP21, LP02 etc.) in addition to the fundamental mode (e.g., LP01). Such higher orders modes (HOMs) tend to degrade the quality of output optical energy provided by the fiber apparatus and hence raise the M2 parameter (lower M2 means better beam quality). In many applications a low M2 is desired. Forestalling the onset of non-linear effects while also maintaining good beam quality can present a challenge to the designer of optical fiber apparatus.\nSome approaches to this challenge are known in the art. For example, U.S. Pat. No. 6,496,301, issued on Dec. 17, 2002 to Koplow, Kliner and Goldberg, teaches bending a multimode fiber having a larger core to substantially attenuate, via increased bend loss, higher order modes such that a fiber amplifier provides gain in substantially a single mode. See also U.S. Pat. No. 7,424,193, issued on Sep. 9, 2008 to Alamantas Galvanauskas, which teaches a composite waveguide having a central core and at least one side core helically wound about the central core and in optical proximity to the central core. According to the '193 patent, higher order modes of the central core selectively couple to the helical side core and experience high loss such that the central core is effectively single-mode.\nExisting techniques, however, are not necessarily entirely satisfactory in all circumstances. Accordingly, it is an object of the present invention to address one or more of the deficiencies or drawbacks of the prior art."}
{"text": "The present invention relates to surgical instruments and, in various arrangements, to surgical stapling and cutting instruments and staple cartridges for use therewith that are designed to staple and cut tissue.\nCorresponding reference characters indicate corresponding parts throughout the several views. The exemplifications set out herein illustrate various embodiments of the invention, in one form, and such exemplifications are not to be construed as limiting the scope of the invention in any manner."}
{"text": "Typically, a pumping unit for an oil well includes a compressor mounted between the walking beam and the Samson post or other stationary part of the pumping unit. A chronic problem associated with walking beam compressors is that rod side loading often occurs which results in excessive wear of the rod seals in the rod seal assembly of the compressor which engage the piston rod as it moves up and down. Various attempts have been made to overcome the problems associated with side loading. For example, Mayland U.S. Pat. No.5,290,156 discloses the use of spherical bearings at each of the connecting ends of the compressor to allow universal movement in response to any lateral displacement or side loading. Studinger U.S. Pat. No. 4,345,744 uses spherical bearings at each of the connecting ends of the compressor to allow non-restricted, omni-directional movement at both connecting ends. McCoy U.S. Pat. No. 4,530,646 has a swivel connector at opposite ends for connecting the compressor to the walking beam and the Samson post to accommodate side loading. McClung U.S. Pat. No. 3,655,301 has knuckle joints and a buckle adjuster to minimize side loading. While these devices are satisfactory for their intended purpose, no attempt has been made to modify the compressor structure to minimize side loads thereon."}
{"text": "The present invention is directed to a support system for electrical cables made out of insulating material for use in high voltage equipment such as liquid filled transformers and reactors. Electrical cables in high voltage equipment such as cables in liquid filled transformers or reactors which connect winding ends or winding taps to tap changers, bushings, or other components like current transformers and the like are subjected to high mechanical forces in the event of a short circuit in service. They are also subjected to over-voltages such as lightning surges which create high electrical stress.\nIt is common practice to custom design the necessary cable clamping fixtures for each individual transformer. The cable clamping fixtures are generally made out of wood, transformer wood or transformer board. It is obvious that when the production is based on the custom fitting of a cable clamping system, the production costs become expensive and lead to an unwarranted lengthening of the production cycle."}
{"text": "Although the present invention can be applied to a wide range of plastics material parts for decorative and trim purposes, for example in the field of automotive engineering, but also in other technical fields, for example when equipping any vehicles, aircraft or buildings, and also in domestic appliances or entertainment devices, the invention and the problem on which it is based are described in detail below with reference to a trim component for a motor vehicle, but without restricting the invention to that effect.\nIn the context of motor vehicles, in particular cars, vendors and users now require not only efficient drive technology, low fuel consumption and the like, but also a high-grade and aesthetic appearance. Decorative or trim components for the interior of motor vehicles are currently often produced from plastics material, at least in part. This allows the cost-effective production of complicated geometries with a simultaneous reduction in weight.\nIn this context, it can happen that a trim component of this type is of a high-gloss design for aesthetic reasons and is to be embellished in particular, predefined regions with a decoration, for example a chrome-look decoration.\nFor this purpose, the trim component can be implemented in two or more separate pieces, one of the pieces being provided with a high-gloss surface and a second of the pieces being provided with a chrome finish or another decoration. For this purpose, the second piece may for example be chromium-plated. However, a two-piece or multi-piece design is associated with high complexity and high costs, especially since the pieces must be produced individually and then assembled.\nThis is a situation which is worth improving."}
{"text": "1. Field of the Invention\nThe invention relates to a master-slave method and device for controlling a thyristor bridge in case of a multiphase current.\n2. Description of the Prior Art\nAs disclosed in particular in the book of H. Lilen \"Thyristor et triac\" published by the Editions Radio, it is known in the art to use a circuit for controlling the time duration of the conduction of a thyristor during an halfwave of the feed voltage in order to control the power passing through a load. Such a circuit has been developed by the Applicant under the name of \"Circuit L1\" and exhibits in particular the advantage to permit a linear control of the feed power through a load.\nIn the prior art, an application of such circuits in a multiphase network and in particular in a three-phase network is disclosed in chapter 13 of the above cited book of H. Lilen, in particular in FIGS. 13-24 of this book which is shown in attachment (FIG. 1). Referring to FIG. 1, three circuits \"L1\" are used for controlling the three thyristors present on the phases U, V and W. Said three circuits \"L1\" are respectively connected for being fed and timed to the phases VU, UW, WV, said connections being made through three coupling transformers (not shown in FIG. 1). Said thyristor control circuit of the prior art presents several drawbacks. The first drawback stands in the necessity to provide for three circuits \"L1\" or analogous or three coupling transformers to permit the control of the power in the three-phase circuit, which increases the costs, the complexity and the bulk of the control device. The second drawback stands in the necessity to suitably connect each circuit \"L1\" to different determined phases of the three-phase feed circuit, whereby the error risk is not negligible. A third drawback stands in the fact that the circuits \"L1\" are circuits which have to be adjusted while inserted in a network; accordingly, three independent adjustments are necessary for balancing the circuit."}
{"text": "The present invention relates to the art of folding paperboard cartons in general and to windowed folding cartons in particular. The windows extend into two or more contiguous walls of the carton and they are covered by transparent thermoplastic sheets. Cartons of this general type are well known to the art and are disclosed, for example, in U.S. Pat. Nos. 4,826,004; 4,846,775; and 4,733,916. It is to improvements in this type of windowed folding carton having an \"around-the-corner\" window that the present invention is directed.\nWhen the flat paperboard blanks from which folding cartons are die cut are formed with large apertures which serve as viewing windows in the finished cartons, those windows extend over to contiguous panels of the blank (eventually walls of an erected carton) and the windows are typically closed off and sealed completely about their peripheries with sheets of transparent plastic sheet. The bonding of the clear plastic sheet to the paperboard about the edges of the window aperture is effected by means of an adhesive which is typically applied in a circumscribing pattern about the window aperture. While cartons erected from this type of windowed blank are very attractive and desirable for the purposes of providing large display areas of packaged goods, the folding and bending of the comparatively stiff plastic sheet along with the underlying paperboard panels often creates tension as well as compression stresses in the joined materials which often results in the tearing of the paperboard along the fold lines at the overlap of the plastic sheet with the paperboard or forms wrinkles in the overlying plastic sheet near the edges of the aperture forming the borders of the window at the area of the fold line or actually destroys the adhesive bonding of the overlying window material to the paperboard itself. Any of these failures, or worse still, several of them simultaneously ruins the appearance and the acceptability of the windowed carton.\nAccordingly, many efforts have been made in composite folding cartons, (those in which window panels are adhered to underlying paperboard blanks in areas which extend over two or more contiguous panels) to eliminate the deleterious aesthetic effects of crinkling, tearing, separation or combinations thereof in the window area. These efforts have concentrated on eliminating or reducing the stressing of the composite blank in the sensitive areas, the zones of the window periphery where the plastic sheet is adhered to the paperboard blank along fold lines between contiguous panels. The present invention provides a solution to this problem which permits the production of folding cartons having a comparatively heavy plastic sheet material closing off display windows formed by apertures extending over two or three contiguous main carton panels in a manner which avoids stress wrinkling of the adhered plastic sheet material, which avoids cracking and/or bursting of the paperboard blank at the score lines where the contiguous panel members are adhered to the plastic sheet, and which avoids delamination, separation or destruction of the sealing adhesive which secures the carton blank to the overlying plastic to the paperboard blank."}
{"text": "In the field of wireless communications, for some applications, it is known to provide a so-called Firmware Over-The-Air (FOTA) system. FOTA systems are used for keeping communications-enabled devices in an updating system updated to their most recent firmware version. The over-the-air update approach enables a device with communication capabilities to be updated using a wireless connection with a server, for example a dedicated firmware updating server. Therefore, the devices can remain in-situ without a need to retrieve or recall them for updates or physical connection to another computing device. However, a number of constraints accompany the over-the-air approach. For example, the device needs to be reachable from a communications perspective, an adequate data transfer rate needs to be available, and the power supply to the device needs to adequate, such as in the case of battery-powered beacons.\nConsequently, updates sent to the device are typically processed in order to account for limited bandwidth in respect of the connection between the updating server and the device, and the fact that wireless reception by the device consumes a relatively large amount of power compared to normal tasks performed by the device. In this regard, for example, the update can be processed by packaging it using a suitable packaging algorithm for the data type involved. This, however, then leads to a requirement that the device receiving the processed update has to be able to reverse the processing of the data in order to use it.\nFurthermore, an update common to all devices using the updating system may not be available. Instead, updates for devices can be required on an individual basis, or groups of more than one device can require respectively different updates to be applied to them. Additionally, in some scenarios, the devices can be accessible for updating purposes by more than one server and so the firmware configuration of a given device is not under the exclusive control of a single updating server. Also, for some devices, it is possible to update the firmware locally, for example by tethering the device to another computing device to apply the update.\nTherefore, whilst an updating server, hereinafter referred to as a “content server”, can be programmed to record and thus track the update status and version of the firmware for a given device, such an approach assumes that the content server is the only server to have access to the device for updating purposes and so there is an assumption that the information held by the content server is always up-to-date. Similarly, whilst the content server can maintain a record of processes supported by each device or group of devices for the purposes of communication of updates from the content server to the device(s), since the content server is not necessarily the only entity to have access to the processes stored by a given device, and it is possible that the device or group of devices are controlled by another server or other servers, or they can be accessed locally for updating purposes. This can lead to changes being made to the device(s) outside the scope of detection of the content server and so the record of processes supported by the content server becomes inaccurate.\nIn order to overcome such difficulties, a known device updating procedure can require the content server to interrogate the device to be updated in order to identify a processing technique, available to the updating server and supported by the device to be updated, which can be used by the content server to communicate, in processed form, the update for the device. However, the known updating procedure assumes that a common processing technique can be found. Additionally, the interrogation of the device requires the device to send data back to the content server in response to a query concerning processing techniques supported by the device, the device thereby consuming power."}
{"text": "One or more embodiments of the subject matter described herein relate to electric motors, such as traction motors of a vehicle. While certain embodiments are described in terms of traction motors of a locomotive or other rail vehicle, the subject matter described herein also may apply to other vehicles and/or other off-highway vehicles (OHV).\nKnown vehicles may include several electric motors, such as three-phase alternating current (AC) traction motors, that receive three-phase AC to power the motors. With respect to each motor, different phases of the current are passed to different conductive coils disposed in a stator of the motor. The current generates a magnetic field in the stator and causes a rotor of the motor to rotate within the stator. The rotor may be coupled with an axle or wheel of the vehicle by one or more gears or other couplings. Rotation of the rotor causes rotation of the axle and wheel to propel the vehicle.\nSeveral mechanical components may be used to enable rotation of the rotor within the stator of the motor. For example, bearings may be disposed between the rotor and the stator to center the rotor in the stator and allow the rotor to rotate at relatively high speeds within the stator. Gears and/or other coupling components may be coupled with the rotor to translate rotation of the rotor to rotation of an axle or wheels. Over time, one or more of the bearings, gears, and/or other coupling components may begin to fail. For example, friction between a bearing and the rotor or stator, friction between gears, and/or friction between two or more other components of the motor may increase as the bearing, gear, or other component begins to mechanically fail. If the motor having the failing bearing, gear, or other component is not identified in time, the failing bearing, gear, or other component may seize or lock up and cause the motor to fail. Once the motor fails, the motor can no longer operate to propel the vehicle.\nSome known systems and methods use additional sensors that are added to the motors of a vehicle to determine if the motor is tending toward failure, such as by a failing bearing, gear, or other component. However, the addition of these sensors can increase the cost and/or maintenance required for the vehicles. Moreover, some of these known sensors may be unable to accurately distinguish between a motor that is tending toward failing and external noise, such as uneven surfaces, tracks, and the like, that the vehicle is traveling over. For example, the uneven surfaces can cause the sensors to detect false positive identifications of a motor tending toward failure.\nA need exists for a system and method for predicting impending mechanical failure of an electric motor that does not significantly add to the cost and/or maintenance of the vehicle and/or can accurately distinguish between failing motors and external noise of the system."}
{"text": "Field of the Invention\nAt least one example in accordance with the present invention relates generally to the control of square wave power supplies.\nDiscussion of Related Art\nThe use of power devices, such as uninterruptible power supplies (UPS), to provide regulated, uninterrupted power for sensitive and/or critical loads, such as computer systems and other data processing systems, is known. Known uninterruptible power supplies include on-line UPS's, off-line UPS's, line interactive UPS's as well as others. On-line UPS's provide conditioned AC power as well as back-up AC power upon interruption of a primary source of AC power. Off-line UPS's typically do not provide conditioning of input AC power, but do provide back-up AC power upon interruption of the primary AC power source. Line interactive UPS's are similar to off-line UPS's in that they switch to battery power when a blackout occurs. Some off-line UPS's include a tap switching voltage regulation circuit to stabilize the output voltage continuously."}
{"text": "1. Field of the Invention:\nThe present invention pertains to brine feeds for use in chlor-alkali cells. More particularly, the present invention pertains to the treatment of brine prior to the electrolytic decomposition thereof. Even more particularly, the present invention pertains to the further treatment of a conventionally treated brine feed prior to the electrolytic decomposition thereof.\n2. Prior Art:\nAs is known to those skilled in the art to which the invention pertains, a sodium chloride brine of relatively high purity is required to operate an electrolytic diaphragm cell for the production of chlorine and sodium hydroxide therefrom. This is because the hardness ions in the brine, e.g. calcium, magnesium and strontium ions, tend to precipitate out of the brine as hydroxides or carbonates during electrolysis. These precipitates form in the diaphragm thereby plugging same. Plugging of the diaphragm reduces cell efficiency and neccessitates the shutting of the cell to either renew or replace the diaphragm.\nThe prior art has appreciated this problem and has sought to alleviate it. For example, in U.S. Pat. Nos. 1,259,683 and 1,308,509 there is taught the addition of sodium chloride to the brine to maintain a high concentration of sodium chloride therewithin in comparison to the hardness ions. Other art teaches the operation of a chlor-alkali cell at currents below the necessary to decompose the calcium salts present in the brine. Still other pertinent prior art is found in U.S. Pat. No. 1,386,244 which teaches a pre-treatment of the feed by adding thereto alkali materials which react with the hardness ions to form precipitates which are then separated out from the brine. The reference specifically teaches the stage-wise precipitation of the hardness ions by the addition of caustic soda to the feed, coupled with the filtration of the solution to provide a purified brine. Because of the cost associated with these prior art processes they are, generally, not commercially employed.\nConventionally, today, the brine is treated by the addition thereto of sodium carbonate and sodium hydroxide to precipitate, simultaneously, magnesium hydroxide and calcium carbonate. The so-treated brine is, then, clarified and filtered. Generally, by so-treating the brine the total hardness ion level (based on calcium) is reduced to about ten parts per million, by weight, based on the weight of the brine. Also, the brine is at a pH level of about eight to ten. However, even this treatment of the brine does not provide sufficient reduction in the hardness ion content to avoid the problems of rapid plugging of the diaphragm.\nThe present invention, on the other hand, reduces the hardness ion level in the brine while being economically feasible and acceptable."}
{"text": "WO 2012/1034141 provided a carbonaceous feed pyrolysis apparatus including two or more hot particle fluidised beds, and one or more positive displacement apparatus for the transfer of hot particles between two or more of the beds. One or more of the fluidised beds may contain a combustion zone. One or more of the fluidised beds may contain a pyrolysis zone. The positive displacement apparatus may be a screw feeder or the like. The apparatus of WO 2012/031414 may include one or more “L” valves between the combustion and pyrolysis zones so that hot particles may flow from the combustion zone or zones to the pyrolysis zone or zones while impeding the flow of gas in the opposite direction.\nThe inventor has now, after extensive research, identified that the apparatus of WO 2012/034141 has practical limitations and that the heat loss makes it impractical to operate without additional fuel and a need was thus identified for a new and efficient apparatus to carry out these and other endothermic processes.\nThe inventor thus now proposes the invention described hereinbelow."}
{"text": "Throwing and catching flying discs is a popular activity among humans as well as between humans and their pets. In use, traditional flying discs can be difficult to catch when in flight at high speeds due to the solid materials from which they are made as well as their unforgiving structure, particularly at the rim of the disc. Traditional flying discs can also be difficult to pick up off the ground depending on the flying disc's orientation as it lays on the surface. For example, when the flying disc is lying “face down” on the ground (grass, concrete, asphalt, etc.) such that the inside of the disc is facing downwards (dome-shape upwards), a user must reach underneath the dome of the traditional flying disc to pick it up. This can be difficult as a user would have to wedge their fingers between the ground and the disc to gain enough leverage to elevate the flying disc. Similarly, dogs attempting to pick up a traditional flying disc lying face down may encounter difficulty getting a firm grasp on the edge of the disc.\nAn improvement is needed over traditional flying discs."}
{"text": "(1) Field of the Invention\nThe present invention relates to a driving circuit for a stepping motor being constituted to perform bipolar driving.\n(2) Description of the Related Art\nA driving circuit of a bipolar system has been used as a driving circuit for a stepping motor.\nA driving circuit of a bipolar system is shown in FIG. 9. As shown in the figure the driving circuit is constituted with a first and a second bridge circuits, 5 and 6, connected between a DC power source 1 and coils, 3 and 4, of a stepping motor 2. The first bridge circuit 5 is formed with a first, a second, a third and a fourth switching elements, Q1, Q2, Q3 and Q4, being composed of transistors, and the second bridge circuit 6 is formed with a fifth, a sixth, a seventh and an eighth switching elements, Q5, Q6, Q7 and Q8 being composed of transistors. In order to feedback the counter electromotive force (accumulated energy) induced in the coils, 3 and 4, to the power source, a first to an eighth diodes, D1 to D8 are connected in antiparallel to the first to the eighth switching elements, Q1 to Q8, respectively.\nAn end of the first coil 3 is connected to the middle point between the first and the second switching elements, Q1 and Q2, and the other end of the first coil 3 is connected to the middle point between the third and the fourth switching elements, Q3 and Q4. An end of the second coil 4 is connected to the middle point between the fifth and the sixth switching elements, Q5 and Q6, and the other end of the second coil 4 is connected to the middle point between the seventh and the eighth switching elements, Q7 and Q8.\nAn end of the DC power source 1 is connected to the junction point 7 (a terminal on the power source side) of the first and the third switching elements, Q1 and Q3, of the first bridge circuit 5, and to the junction point 8 (a terminal on the power source side) of the fifth and the seventh switching elements, Q5 and Q7, of the second bridge circuit 6, and the other end of the DC power source 1 is connected to the junction (ground) of the second and the fourth switching elements, Q2 and Q4, of the first bridge circuit 5, and to the junction point (ground) of the sixth and the eighth switching elements, Q6 and Q8, of the second bridge circuit 6.\nAn exciting signal generator 9 is a well-known circuit which transmits exciting signals shown in FIGS. 11(A), 11(B), 11(C) and 11(D), to lines, 10, 11, 12 and 13, to drive a stepping motor 2 with a predetermined exciting system (for example, a bipolar exciting system). The exciting signal line 10 is connected to the control terminals (bases) of the first and the fourth switching elements, Q1 and Q4, the exciting signal line 11 is connected to the control terminals of the second and the third switching elements, Q2 and Q3, the exciting line 12 is connected to the control terminals of the fifth and the eighth switching elements, Q5 and Q8, and the exciting signal line 13 is connected to the control terminals of the sixth and the seventh switching elements, Q6 and Q7. A resistor is connected to each of the base lines of switching elements, Q1 to Q8.\nThe principle of the constitution of the stepping motor 2 is as shown in FIG. 10; for example, it comprises a stator core 18 having a first, a second, a third and a fourth magnetic poles, 14, 15, 16 and 17 being disposed at intervals of 90 degrees and a rotor 19 being composed of a permanent magnet having an N pole and an S pole at intervals of 180 degrees. The first coil 3 is wound on the first and the third magnetic poles, 14 and 16, and the second coil 4 is wound on the second and the fourth magnetic poles, 15 and 17.\nWhen the stepping motor 2 is driven by a bipolar two-phase exciting system, the first to the eighth switching elements, Q1 to Q8, are ON/OFF controlled by the exciting signal shown in FIGS. 11(A), 11(B), 11(C) and 11(D). In the period of time, t1 to t3, shown in FIG. 11(A), the first and the fourth switching elements, Q1 and Q4, are made ON simultaneously. Thereby, a current in a first direction is made to flow through the first coil 3 in a circuit constituted with the DC power source 1, the first switching element Q1, the first coil 3 and the fourth switching element Q4.\nIn the period of time, t3 to t5, shown in FIG. 11(B), the second and the third switching elements, Q2 and Q3, are made ON, and a current in a second direction is made to flow through the first coil 3 in a circuit constituted with the DC power source 1, the third switching element Q3, the first coil 3 and the second switching element Q2.\nIn the period of time, t2 to t4, shown in FIG. 11(C), the fifth and the eighth switching elements, Q5 and Q8, are made ON, and a current in the first direction is made to flow through the second coil 4 in a circuit constituted with the DC power source 1, the fifth switching element Q5, the second coil 4 and the eighth switching element Q6.\nIn the period of time, t4 to t6, shown in FIG. 11(D), the sixth and the seventh switching elements, Q6 and Q7, are made ON, and a current in the second direction is made to flow in the second coil 4 in a circuit constituted with the DC power source 1, the seventh switching element Q7, the second coil 4 and the sixth switching element Q6.\nWhen any one of the switching elements, Q1 to Q8, is changed from ON to OFF, a counter electromotive force is induced in the first coil 3 or the second coil 4. When the second and the third switching elements, Q2 and Q3, are changed from ON to OFF, and the first and the fourth switching elements, Q1 and Q4, are changed from OFF to ON, based on the counter electromotive force induced in the first coil 3, a current is made to flow in a circuit constituted with the first coil 3, the first diode D1, the DC power source 1 and the fourth diode D4, and the accumulated energy in the first coil 3 is fed back to the DC power source 1. At this time, since the first and the fourth diodes, D1 and D4, are made ON, the application of a high voltage based on the counter electromotive force to the first and the fourth switching elements, Q1 and Q4, is prevented; thus the first and the fourth switching elements, Q1 and Q4, can be protected.\nWhen the switching elements, Q1, and Q4 to Q8, other than the above-mentioned second and the third switching elements, Q2 and Q3, are changed from ON to OFF, a similar action to the above is performed.\nWhen a current based on the counter electromotive force is flowing through the first and the fourth diodes, D1 and D4, for example, even if exciting signals are given to the first and the fourth switching elements Q1 and Q4, these switching elements, Q1 and Q4, are not made ON. Therefore, when the period of time in which a current based on the counter electromotive force flows (feedback current period) is long, since an actual exciting period is disposed after the feedback current period, it becomes necessary to lengthen the period of the exciting signal; thereby, there occurs a problem that the stepping motor 2 cannot be operated at a high rotational speed.\nA method described in the following is well-known to the public as a method for solving the problem: in the method, in order to make the feedback current period short and to increase the pullout torque of the stepping motor 2, as shown in FIG. 12, a first and a second capacitors, 20 and 21, are connected in parallel to respective bridge circuits, 5 and 6.\nIn FIG. 12, the first capacitor 20 is connected between the junction point 7 of the first bridge circuit 5 and the ground, and a diode 22 for preventing feedback is connected between the junction point 7 and the DC power source 1. The second capacitor 21 is connected between the junction point 8 of the second bridge circuit 6 and the ground, and a diode 23 for preventing feedback is connected between the junction point 8 and the DC power source 1. Except the above-mentioned points, the circuit is constituted in the same manner as that shown in FIG. 9.\nIn the circuit shown in FIG. 12, owing to the capacitors, 20 and 21, for example, as shown in FIGS. 13(A) and 13(B), at the time of t1, the exciting signal is given to the first and the fourth switching elements, Q1 and Q4, and when the second and the third switching elements, Q2 and Q3, are made OFF, a current based on the counter electromotive force induced in the first coil 3 is made to flow through a circuit constituted with the first coil 3, the first diode D1, the first capacitor 20 and the fourth diode D4. Thereby, the accumulated energy in the first coil 3 is absorbed by the capacitor 20. The period of time in which a current flows based on the counter electromotive force can be adjusted by the capacity of the capacitor 20, so that it is made possible to operate the stepping motor 2 at a high rotational speed by making the feedback period shorter than that in the circuit described in FIG. 9. The energy accumulated in the capacitor 20 functions as a power source to make a current flow in the first direction in the first coil 3 when the first and the fourth switching elements, Q1 and Q4, are made ON; therefore, it is made possible to increase the pullout torque of the stepping motor 2.\nIn the circuit shown in FIG. 12, however, when the charged voltage of the capacitors, 20 and 21, is made higher than the breakdown voltage of the switching elements, Q1 to Q8, by the dispersion of circuit constants such as the first and the second coils, 3 and 4, or the capacitors, 20 and 21, the breakdown of the switching elements may occur in the OFF period and in some case they may be destroyed. In FIG. 13(C), the charged voltage Vc of the capacitor 20 is shown, and at a time t1, when the switching elements, Q2 and Q3, are turned OFF (change from ON to OFF), it is charged up to a voltage Vsp which is higher than the source voltage Vs, and if the Vsp becomes higher than the resisting voltage of the second switching element Q2 or of the third switching element Q3, there used to be a fear that the breakdown of a switching element can occur."}
{"text": "1. Field\nThe disclosure relates generally to extending operational capabilities of mobile devices, and more specifically to extending file download capabilities of mobile devices having diminished memory capacity.\n2. Description of the Related Art\nIn this day and age, with the increasing need to deliver multimedia content to the widest variety of users and devices, the lack of sufficient device memory space of where to store the content, and then use it, is being encountered more frequently.\nThe Session Initiation Protocol (SIP) is a signaling protocol defined by an Internet Engineering Task Force (IETF) and widely used for controlling multimedia communication sessions such as voice and video calls over Internet Protocol (IP). The protocol can be used for creating, modifying and terminating two-party (unicast) or multiparty (multicast) sessions consisting of one or several media streams. The modification can involve changing addresses or ports, inviting more participants, and adding or deleting media streams. Other feasible application examples include video conferencing, streaming multimedia distribution, instant messaging, presence information, file transfer and online games. The new SIP enabled devices allow users to be reached by every content type (voice, video, images, and data), including mixed content, and signal their presence or availability by publishing their state to a presence server.\nThe new generation mobile phones are usually sold with a small memory card, which afterwards, can be extended with greater and more expensive cards. And even when extended, for example to 512 MB, memory can be partially occupied when there is the need of downloading a video-conference or other big multimedia files, so that there is not enough free space and not enough time to select data to be deleted to free memory.\nIt often happens in specific conditions or during specific events like a business meeting in which multiple devices are available but each one of them cannot work with other ones in order to have a cooperative download of a multimedia file that may provide an optimized and cheaper download method.\nThe following described techniques address these issues."}
{"text": "Sputtering is a technique used to produce a metallic layer in various manufacturing processes used in the semiconductor and the photoelectric industries. The properties of films formed during sputtering are related to the properties of the sputtering target itself, such as the size of the respective crystal grain and the formation of secondary phase with distribution characteristics. It is desirable to produce a sputter target that will provide film uniformity, minimal particle generation during sputtering, and the desired electrical properties.\nVarious sputtering techniques are used in order to effect the deposition of a film over the surface of a substrate. Deposited metal films, such as metal films on a flat panel display device, can be formed by a magnetron sputtering apparatus or other sputtering techniques. The magnetron sputtering apparatus induces plasma ions of a gas to bombard a target, causing surface atoms of the target material to be ejected and deposited as a film or layer on the surface of a substrate. Conventionally, a sputtering source in the form of a planar disc or rectangle is used as the target, and ejected atoms travel along a line-of-sight trajectory to deposit on top of a wafer whose deposition face is parallel to the erosion face of the target. Tubular-shaped sputtering targets can also be used, as described in co-pending application Ser. No. 10/931,203.\nSputtering targets may be desired which comprise materials or combinations of materials that cannot be made by conventional means such as rolling. In such cases, targets are made by hot isostatic pressing (HIP) powders. Ideally, the target is made in a single step. However, physical limitations of powder packing density and size of HIP equipment make it necessary to join smaller segments in order to produce large sputtering targets. For single-phase targets, conventional processing such as welding may be used; for multi-phase materials or where alloy formation is to be avoided for any reason, solid-state edge to edge bonding is preferred.\nInterconnects in semiconductors and TFT-LCDs are evolving from aluminum and toward copper, thus new diffusion barriers are needed. Titanium provides excellent adhesion properties while the molybdenum contributes its dense barrier stability. Integrated circuits (for semiconductors and flat panel displays) use Mo—Ti as an underlayer or capping layer for aluminum, copper, and aluminum alloys to minimize hillocks formation, to control the reflectivity and provide protection from chemical attack during photolithography.\nU.S. Pat. No. 5,234,487 describes methods of producing tungsten-titanium sputter targets with little or no β(Ti, W) alloy phase. U.S. Pat. No. 5,896,553 describes a titanium-tungsten sputter target which is substantially all single phase β(Ti, W). Neither patent discloses the use of other materials as substitutes for tungsten.\nU.S. Patent Application Publication 20050189401 discloses a method of making a large Mo billet or bar for a sputtering target wherein two or more bodies comprising Mo are placed adjacent one another (e.g. stacked one on the other) with Mo powder metal present at gaps or joints between the adjacent bodies. The adjacent bodies are hot isostatically pressed to form a diffusion bond at each of the metal-to-Mo powder layer-to-metal joint between adjacent bodies to form a billet or bar that can be machined or otherwise formed to provide a large sputtering target. This patent publication discloses bonding of major side surfaces, not edge-to-edge bonding of plates."}
{"text": "1. Field of the Invention\nThis invention relates to a disposable-type indwelling needle set provided with a needle tip composed of a soft synthetic resin.\n2. Description of the Related Art\nThere is a conventionally known indwelling needle set comprising a metallic inner needle, which is provided with a flexible outer needle composed of a synthetic resin, outside the inner needle. A patient is pierced with the tip of the inner needle to insert the tip of the outer needle in a blood vessel together with the inner needle. Then, only the inner needle is removed to connect a connecting tube to the base of the outer needle whereby treatments such as drip infusion or the like are carried out.\nHowever, complicated and careful operations are required for preventing infection and the like to carry out the sequence using the conventional indwelling needle set, in which only the inner needle is removed under such conditions that a blood vessel is pierced with the outer needle and then the connecting tube is connected to the base of the outer needle."}
{"text": "High purity and highly corrosive liquids are often utilized in industries such as semiconductor processing and chemical manufacturing. Due to quality and safety concerns, these liquids generally must be contained in high integrity containers and dispensed using highly reliable dispensing systems.\nConventionally, the containers are drums formed of high purity polyethylene, such as those described in U.S. Pat. No. 6,045,000. Examples of dispensing systems suitable for use with these drums are described in U.S. Pat. Nos. 4,699,298, 5,108,015, 5,957,328, and 5,526,956. The connections, tubing, and fittings for handling these fluids are often formed from inert materials, such as various fluoropolymers. The tubing is commonly formed of Per Fluor Alkoxy (“PFA”), while the fittings and valve components can be formed of components such as PFA, Poly Tetra Fluoro Ethylene (“PTFE”), and other various fluoropolymers.\nConventional drums often include a bung opening and a drum insert connected to an upward and concentric nipple on the bung. A dispense head can be operably coupled with the drum insert and can include a main fluid passageway that connects to and sealingly engages the bung.\nIn many applications, the fluid that is dispensed from the drum can later be recirculated to the drum. This has been accomplished by using a recirculation passageway in the dispense head. However, these recirculation passageways can be uncontrolled and thus can remain in an open position when the dispense head is not connected to the drum, which can create a potential for spillage.\nAlso, recirculated fluid can create foam that can cause problems, such as foam entry into the vent lines, potential clogging, contamination of normally dry areas of the dispense head and gas refill line, and difficulty in emptying a drum with low levels of fluid. The foam can be created when fluid is recirculated into the dispense head and further when recirculated fluid is dispensed into the drum.\nWhen using dispense heads to draw high purity, highly corrosive, and highly caustic liquids from drums, it can be important not to connect a dispense head for equipment requiring one type of liquid to a drum containing an incorrect or incompatible different liquid. To inhibit this, coded indexing sections and recesses have been provided in the dispense head and drum insert to enable assembly of the drum insert to the dispense head only when the correct coding exists. However, connecting a coded dispense head to a coded drum insert can require a great amount of force while turning the coding and holding down the dispense head. This process can involve a two-step operation of holding down the dispense head with one hand while rotating the code key with the other."}
{"text": "The invention relates to a spout on filling apparatus for packing liquids, having a spout plate which has a plurality of bores disposed closely side by side for the passage of the liquid into a container or the like disposed underneath the spout.\nThe spouts on packing apparatus of the kind herein contemplated are commonly operated cyclically, so that the material being packed, such as milk or fruit juice, can be poured intermittently into the containers moved successively and cyclically under the spout. Operation at maximum speed in the sense of short cycling times requires that the stream emerging from the spout be made as large as possible in relation to the container cross section, which in turn requires that the filling stream be given a precise shape with a smooth outside surface. The intermittent manner of operation furthermore necessitates measures for preventing the liquid from dripping between fill cycles, in order to prevent contamination of the apparatus and trouble in closing the containers such as might be caused by wetting them with the liquid in the area of the closure. For the achievement of this purpose in the spouts it is known to dispose strainer or filter disks in the spouts, or else spout plates having bores disposed closely adjacent one another. The invention relates to this last-described type of spouts. This type has the advantage over the use of strainer and filter disks, which cannot be used with pulpy liquids, of a wider range of applications. However, the known spouts of this kind are not free from disadvantages, either. For example, the individual bores in the spout plate result in a great number of individual jets having a relatively large total surface area, resulting in severe frothing. The froth on the surface of the material in the filled container is disadvantageous for the reason that it can again cause wetting in the container areas which have to be heated afterward for closing by heat sealing or the like. The larger surface area of the many individual jets, however, also causes fine air bubbles to be entrained in the liquid, which initially do not immediately manifest themselves in froth, but do raise the liquid level and thus also cause difficulty in closing.\nIt is already known to dispose the bores in the spout plate at a low angle to the axis of the spout, or to set the container at an angle underneath the spout, so that the liquid jets emerging from the bores impinge upon the inside wall of the container at an angle, and the liquid flows down this wall in the form of a film (cf. European patent application No. 00 13132). In this manner a certain improvement can be achieved as regards foaming and the entrainment of air into the liquid, but it does not go all the way, for the problem still persists of the dripping of the liquid between the individual container filling cycles."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display device used for a display part of an electronic equipment, and particularly to a liquid crystal display device including a vertical alignment type liquid crystal having a negative dielectric anisotropy.\n2. Description of the Related Art\nA vertical alignment (VA) mode liquid crystal display device has features of a high contrast ratio and a high speed response characteristic, and in recent years, the VA mode liquid crystal display device has been actively developed. Especially, an MVA (Multi-domain Vertical Alignment) mode liquid crystal display device has features of a wide viewing angle, a high contrast ratio and a high speed response characteristic, and receives attention as a display system most suitable for a flat panel display for a TV receiver. FIG. 7 shows a sectional structure of one pixel of the MVA mode liquid crystal display device. As shown in FIG. 7, the MVA mode liquid crystal display device includes a glass substrate 103 of a TFT substrate 102, an opposite side glass substrate 105 of an opposite substrate 104 disposed to be opposite to the TFT substrate 102, and a liquid crystal 112 sealed between both the substrates 102 and 104.\nAn insulating layer 106 is formed on the glass substrate 103. Plural gate bus lines and plural drain bus lines (neither of them are shown) intersecting each other through the insulating layer 106 are formed on the glass substrate 103. A not-shown thin film transistor (TFT) is formed at each of intersection parts of the gate bus lines and the drain bus lines. The gate bus lines and the drain bus lines are insulated from each other through the insulating layer 106. Besides, the insulating film 106 functions as a gate insulating film of the TFT. A pixel electrode 110 made of indium tin oxide (ITO) and patterned into a specified shape is formed on a final protection film 108. The pixel electrode 110 is connected to a source electrode of the TFT through a not-shown contact hole formed in the final protection film 108. A pixel electrode slit 116 as an alignment regulating structure for regulating an alignment direction of a liquid crystal molecule 120 is formed in the pixel electrode 110. A vertical alignment film 114 for vertically aligning the liquid crystal molecule 120 is formed on the whole surface of the pixel electrode 110 and the pixel electrode slit 116.\nOn the other hand, a not-shown color filter (CF) layer is formed on the opposite side glass substrate 105. An opposite electrode 124 made of ITO is formed on the CF layer and the whole surface of the substrate. A linear protrusion 118 as an alignment regulating structure protruding on the opposite electrode 124 is formed on the opposite side glass substrate 105. Similarly to the pixel electrode slit 116, the linear protrusion 118 is formed in order to regulate the alignment direction of the liquid crystal molecule 120. A vertical alignment film 122 is formed on the whole surface of the opposite electrode 124 and covers the linear protrusion 118. In the MVA mode liquid crystal display device, the linear protrusion 118 and the pixel electrode slit 116 are provided in the pixel, so that the alignment control of the liquid crystal 112 and multi-domain formation are realized.\n[Patent document 1] JP-A-2001-183637\n[Patent document 2] JP-A-8-122753\nAs a method of keeping the thickness (cell gap) of the layer of the liquid crystal 112 at a desired length, there is used a method of scattering spherical ball spacers 136 each having a diameter equal to the desired cell gap into the layer of the liquid crystal 112. The ball spacers 136 are made of plastic material or glass material. The ball spacers 136 are scattered on the TFT substrate 102 or the opposite substrate 104, and the TFT substrate 102 and the opposite substrate 104 are attached to each other through a seal material formed into a frame shape, so that the liquid crystal 112 is sealed between both the substrates 102 and 104.\nFor example, when the ball spacers 136 are scattered on the opposite substrate 104, there is a case where the ball spacer 136 is disposed on the linear protrusion 118. The cell gap at the position where the linear protrusion 118 is formed is narrow as compared with the other position. Thus, as shown in FIG. 7, the cell gap at the linear protrusion 118 is narrower than the diameter of the ball spacer 136. Since the ball spacer 136 is made of relatively hard material as compared with the vertical alignment films 114 and 122, when the opposite substrate 104 and the TFT substrate 102 are attached to each other, there is a case where the vertical alignment films 114 and 122 on the linear protrusion 118 and on the opposite side thereof are damaged by the pressure of the ball spacer 136.\nAt this state, the alignment control of the liquid crystal molecule 120 at the position becomes impossible, light leak occurs at the linear protrusion 118 and in the vicinity thereof, and a poor display occurs in the MVA mode liquid crystal display device. In general, in order to realize excellent alignment of the liquid crystal molecule 120, it is necessary that the width of the linear protrusion 118 is about 10 μm. This corresponds to several percents of the area of the pixel region. For example, in a 15-inch MVA mode liquid crystal display device with a resolution of XGA (extended Graphics Array), in the case where the linear protrusion 118 is formed to have a width of 15 μm, and 100 ball spacers per 1 mm2 are scattered, there is a case where about 7 ball spacers per 1 mm2 are scattered on the linear protrusion 118. When the vertical alignment films 114 and 122 are damaged by the ball spacer 136, the light leak occurring on the display screen becomes very noticeable."}
{"text": "1. Field of the Invention\nThis invention relates to zoom lenses and, more particularly, to zoom lenses of compact form, which nevertheless get a high zoom ratio.\n2. Description of the Related Art\nRecent minimization of the bulk and size of leaf shutter type cameras or the like has led to a demand for shortening the total length of the compact zoom lens.\nIn Japanese Laid-Open Patent Application No. Sho 57-201213, there is proposed a compact zoom lens comprising two lens units of plus-minus refractive power from the object side with the air separation between these lens units being varied to effect zooming.\nAlso, the use of three lens units of minus-plus-minus refractive power in the zoom lens, which are movable for zooming, is proposed in, for example, Japanese Laid-Open Patent Applications No. Sho 63-271214, No. Sho 64-72114, etc. In these publications, the relation in which the first, second and third lens units move during zooming, and the construction and arrangement of the constituent parts of the three lens units are specified to thereby well correct any variation of aberrations with zooming. Thus, a relatively compact lens system having a high optical performance at a zoom ratio of 2 to 3 with 8 or 9 lens elements as a whole is achieved.\nAn example of adaptation of the zoom lens to the leaf shutter camera by extremely reducing the number of lens elements is proposed in Japanese Laid-Open Patent Application No. Sho 63-276013.\nIn this Japanese Laid-Open Patent Application No. Sho 63-276013, the front lens unit comprises negative and positive, refractive power or two, lenses, while the rear lens unit comprises positive and negative refractive power lenses. Thus, a zoom lens composed of four lens elements in total is achieved. However, because the zoom lens uses a refractive power distribution type lens (GRIN lens), the difficulty in manufacturing the zoom lens is extremely high, causing the amount of lenses produced through mass production to become low.\nFurther, Japanese Laid-Open Patent Application No. Sho 64-88512 (U.S. Pat. No. 4,756,609) discloses a zoom lens having four lens units of a minus-plus-plus-minus refractive power arrangement, of which the first and second lens units each comprise one lens element, the third lens unit is a cemented lens, and the fourth lens unit comprises three lens elements of plus-minus-minus refractive power.\nAlso, in Japanese Laid-Open Patent Application No. Sho 62-235916, there is disclosed a zoom lens having four lens units of a minus-plus-plus-minus refractive power arrangement, wherein the first lens unit is constructed with two elements of plus-minus refractive power, the second lens unit with three elements of plus-plus-minus refractive power, the third lens unit with one element, and the fourth lens unit with two elements of minus-plus refractive power."}
{"text": "This invention relates to a leak detection apparatus and method for detecting leaks from a high-pressure or low pressure vessel used to transport and/or store a substance.\nHazardous and flammable substances are often stored and transported in high-pressure vessels or other types of sealed vessels, such as railroad tankcars, semi-tractor trailer tanks, or other vehicles capable of hauling or holding a high-pressure vessel. Often these vessels can develop leaks around the valves and/or other access means used to fill, empty, and/or otherwise access the vessel, and/or pressure relief valves, due to neglect, environmental factors, or even normal use. Intermittent, temporary leaking can also occur during transport due to slamming of a manway lid of the vessel, rough road or rail conditions, and/or the impact of the vessel during humping of the railcars or hitching the semi-tractor trailer.\nLeaking can be expensive to the shipments owner and dangerous to those working near the vessel. Current commercial leak monitoring systems are brought to the vessel to detect leaks and are used when the vessel is stationary, including during loading, unloading, and maintenance conditions. There is a need for a leak detection apparatus useful for detecting leaks of a vessel, such as a railroad tankcar, while in motion or while stationary. There is a need for a detection apparatus that can monitor shipments while the vessel is in motion or stationary and to transmit the detected information to a remote unit. There is a need for a system of detection apparatuses that can monitor information on a plurality of shipment vessels while stationary or in motion and transmit that information through wireless means to a central monitoring unit.\nIt is one object of this invention to provide a system of chemical sensors used to detect leaks of hazardous or flammable substances from vessels during transit, docking, and storage.\nIt is another object of this invention to provide a leak detection apparatus to protect the health and safety of workers and the environment.\nIt is another object of this invention to provide a system for monitoring safety, security, and inventory management of a shipment for shippers, operators, customers, or others involved in the mobile shipment of chemicals and hazardous materials.\nIt is another object of this invention to use at least one chemical sensor and/or other sensor type in combination with wireless technologies, such as a cell phone, a pager, wireless Internet, and/or other satellite technology, and/or Internet technology to transmit the status of a shipment, both while the shipment vessel is in motion and stationary, to a user designated site or a central site. The central site can distribute the status information via dedicated line or via a means such as the Internet, a pager, a cell phone, and/or other electrical or digital communication devices to persons having an interest in the shipment information.\nThe above objectives are accomplished, in one embodiment of this invention, by a leak detection apparatus having a first chemical sensor in communication with a first chamber of a shroud of a vessel and a second chemical sensor in communication with an ambient environment surrounding the shroud. Upon exposure to a predetermined substance, the first and second chemical sensors have a first and second chemical sensor output, respectively. A monitoring circuit can monitor the first chemical sensor output and the second chemical sensor output and detect a first increase of the first chemical sensor output and/or a second increase of the second chemical sensor output. The monitoring circuit signals an information unit in communication with the monitoring circuit. The signal can depend on the first increase and/or the second increase. The signal can also depend on no increase in the first and second chemical sensor outputs, thereby signaling favorable conditions. By comparing the first and second chemical outputs, the leak detection apparatus of this invention can be used effectively while the vessel is in motion. Also, the use of two sensor signals can reduce interference and false alarms because physical and chemical interferences like rapid temperature changes or high pollution levels can be determined by use of the second sensor.\nIn one embodiment, the shroud provides a barrier between the first and second chemical sensors resulting in detectable differences between the first and the second chemical sensor outputs. If a leak occurs in a chamber of the shroud the first chemical sensor output may have a greater increase and/or rate of increase than the second chemical sensor output. If a leak occurs in the ambient environment surrounding the shroud the second chemical sensor output will have a greater increase and/or rate of increase than the first chemical sensor output.\nThe information unit can be an alarm and/or other information reporting means that can inform of hazardous leaks or other information collected by the chemical and/or physical sensors and the monitoring circuit. The information unit can be a local information unit in proximity to the shroud and/or a remote information unit in communication with a transmitter and/or the monitoring circuit. The information unit can emit a signal showing the status of the sensors to various output devices. When the information unit is remote, the monitoring circuit emits an output signal to a transmitter which emits a status signal to a receiver in communication with the remote information unit. The remote information unit can receive signals from more than one transmitter, thereby monitoring more than one vessel. Shippers, operators, and customers can have access to the local and/or remote information units through various output devices, such as computers, cell phones, pagers, and personal digital assistants, to receive information such as the vessel location and vessel integrity associated with their shipments.\nIn another embodiment of this invention, one sensor is sufficient to detect the leak in the area near the vessel valves if the sensor can be adequately compensated for temperature changes, environmental effects, and interferences.\nIn another embodiment of this invention, at least one physical sensor is in communication with the monitoring circuit. The physical sensor can monitor temperature, pressure, movement and/or sounds around the vessel, and/or determine the position of the vessel by means such as satellite positioning technology. The physical sensor is in communication with the monitoring circuit and therefore can also be in communication with the remote information unit and provide information to interested parties."}
{"text": "1. Field of the Invention\nThis invention relates to the field of software, and more particularly to software based simulation tools that use process simulation techniques to assess the performance reliability of products that are made up of multiple components, each component having a reliability distribution.\n2. Description of the Related Art\nIn general, products that have multiple components fail in their intended purpose when a component necessary for the product to perform the function for that purpose fails. By way of illustration, when an automobile, which has many components, has the mission of transporting a commuter to and from work suffers a failure of an electrical component that is necessary for the engine to operate, the automobile is stalled and unable to fulfill its mission. The auto might also serve other purposes, for example transporting people at night to social events. In that case, the failure or burnout of the headlamps of the car makes it unsafe to use at night, and so it cannot fulfill a subsidiary, but nonetheless important, intended purpose.\nIt would be useful to the user of a product to know, or to have some idea, when a component will fail that would create an inconvenience or that would impair the usefulness of the product. In markets where many manufacturers compete to make sell a consumer product, the manufacturer conducts tests to determine component reliability and sets specifications for manufacture that often reflect some reliability goals. Based on an understanding of the reliability of components, gathered from actual failures of components in use, or from test data, or both, a manufacturer can develop a “limited warranty” program. Such a program offers some compensation (cash, or component repair or replacement) to a purchaser whose product has a component failure before the component has reached a certain age or has been in use for a certain time period. Thus, for example, auto manufacturers offer “limited drivetrain warranties” based on the expected trouble free life (measured in mileage used, or age) of the drivetrain they supply. In doing so, manufacturers bear in mind that some few might fail before the warranty period, and that compensation then must be paid. But, manufacturers usually have carried out a cost benefit analysis, using statistical and financial techniques, from which they can determine that the number of failures in the warranty period and the cost of compensation is significantly outweighed by increased sales of their autos (attributable to the limited warranty), and the increased profits from these sales.\nOf course, the consumer product model explained above does not apply to all kinds of products. For example, in the market for large commercial aircraft that carry hundreds of passengers, mission failure cannot be tolerated as readily as in the automotive industry. Accordingly, other techniques are used. In general, these are based on testing of components to determine rates of wear and component life, and observing component wear and life in service. Based on these, regular scheduled maintenance procedures are applied to replace each component on a routine basis before it reaches a point of failure. This technique is also used in other capital industries where the continuous reliable operation of large machines is necessary or critical.\nThere are also other areas where reliability of mission performance poses challenges. For example, some products are manufactured in bulk, and stored for a long period before their intended use. These products each a have a shelf life based on the rate of degradation of their components under storage conditions. For example, a pharmaceutical company may manufacture a medicine that contains several chemical compounds in batches at intervals and supply pharmacies from its inventory. Under these conditions, it is necessary to ensure that the beneficial chemical composition of the medicine not degrade during the storage period. Generally, pharmaceutical companies mark their products with an “expiration date” that is a date after which the company cannot guarantee that the medicine retains its potency. The pharmaceutical industry is not the only one faced with these issues of batch production and storage. Some industries have an added complication that their products have chemical, electronic and mechanical components or subsystems that are each prone to degradation with age, and mission performance reliability must be estimated in some manner."}
{"text": "Display screens of electronic devices such as a mobile phone, personal data assistants (PDAs), personal computers, game consoles and so on, are developing towards big sizes. So, the electronic device generally includes a user interface with a big size."}
{"text": "Inkjet printing technology has become widely known and is often implemented using a thermal ink propulsion mechanism. Such technology forms characters and images on a medium, such as paper, by expelling drops of ink in a controlled fashion so that the drops land on the medium. In its simplest form, the printer can be understood as a working combination of a mechanism for moving and placing the medium in a position such that the ink drops can be placed on the medium, a print cartridge which controls the flow of ink and selectively expels drops of ink to the medium, and appropriate controlling hardware and software to position the medium and expel droplets so that a desired graphic is formed on the medium.\nA conventional print cartridge for an inkjet printer comprises an ink containment device and an ink ejection assembly which selectively heats and expels ink drops in a controlled fashion. An alternative ink ejection apparatus for other than thermal inkjet printers uses mechanical propulsion techniques such as piezo-electric drivers to impart momentum to the ink. For thermal inkjet printers, the ink ejection assembly typically includes a laminate structure comprising a semiconductor or insulator substrate capped with a plate or surface permeated by orifices or nozzles arranged in a pattern beneficial to the deposition of ink on a medium. In thermal inkjet printers, the heating and expulsion mechanism consists of a plurality of heater resistors formed on the semiconductor, each heater resistor associated with one of the nozzles in the nozzle plate. Each of the heater resistors is coupled to the controlling hardware and software of the printer such that each of the resistors may be independently energized to quickly vaporize a portion of ink and expel a drop of ink from the nozzle. Ink is stored in an ink reservoir in the print cartridge, in most implementations, and is typically loaded into a firing chamber which is disposed about the heater resistor. The pressure created following the ink expulsion causes ink to be removed from the ink reservoir through ink manifolds and ink conduits and into the ink firing chambers.\nCurrently, the most widely used thermal inkjet printers utilize a relatively small printhead (approximately 5mmxc3x9710 mm) that is mounted on a protruding xe2x80x9cnosexe2x80x9d of the ink cartridge, thereby being positioned in close proximity to the medium to be printed upon. Such an arrangement is shown diagrammatically in FIG. 1A. Here, a cross section of the nose of the cartridge body 101 shows a printhead 103 affixed to the body 101 in a well 105 having a depth that is approximately equal to the thickness of the printhead 103. A flexible electronics circuit commonly known as a tape automated bonding (TAB) circuit 107 is conventionally affixed to that portion of the body 101 which is directly opposite the medium to be printed upon. The TAB circuit 107 is typically wrapped around one of the other surfaces of the nose so that an electrical connection may be made to the other controlling portions of the printer. This is illustrated diagrammatically with point contact 109. Electrical connection between the TAB circuit 107 and the printhead 103 is typically made using a beam lead technology (as illustrated with beam lead 111) but may also be accomplished with wire bonding techniques. In order to prevent electrical shorts due to ink and to provide mechanical protection, an encapsulating substance (not shown) is typically placed over the wire bonding or beam lead area. Ink is provided to the printhead from a reservoir of ink by ink ducts 113 within the print cartridge.\nAn alternative arrangement which has come into recent use is illustrated in FIG. 1B. This arrangement utilizes the TAB circuit 107xe2x80x2 as an nozzle plate as well as it conventional electrical interconnect purpose. The nozzles are laser ablated or otherwise created in the flexible polymer TAB circuit tape and the TAB circuit is stretched over the remainder of the printhead 103xe2x80x2 in such a way that ink is ejected from the TAB circuit nozzles. As described relative to FIG. 1A, the TAB circuit is wrapped around the side of the body 101 and electrical connection is made by way of a contact 109. Ink is brought to the printhead 103xe2x80x2 by ink ducts 113 in the print cartridge.\nThese techniques have demonstrably proven their value in many successful products. The extension of these technologies to a print cartridge which is expected to print a wide print swath, perhaps the width of an entire medium page, exposes problems which require solutions before these technologies can be used for a wide area print cartridge.\nIf one desires to create a print cartridge which is capable of printing a wide swath in a single pass (for example a 5 cm to 20 +cm swath), or as much as a 20 cm swath, a large number of printing nozzles must be positioned near the medium to be printed upon. It is conceivable that a printhead having thousands of nozzles could be so positioned, but the reliable fabrication of such a large number of nozzles in a single printhead is not economically feasible. Accordingly, a smaller number of nozzles in an ejecting mechanism can be economically produced and a plurality of these ejection mechanisms can be ganged together in a single printhead to effectively obtain a large number of nozzles. Coordinating such a plurality of ink ejection devices requires a high degree of registration of nozzles of separate ink ejectors as well as a high degree of planarity between the ejecting surfaces of the ink ejectors so that the drops ejected from one ink ejector are properly placed on the medium relative to the drops ejected from another ink ejector.\nAccordingly, it is desirable to create a printhead employing a plurality of ink ejection apparatus in a wide array configuration. Such an array should provide a reliable fluidic and electrical interconnect between the print head and the plurality of ink ejection apparatus as well as a high degree of registration between the ink ejection apparatus with a corresponding reduction of misplaced ink drops.\nThe present invention encompasses a printhead for an inkjet printer which employs an essentially rigid substrate having a first surface and a second surface. A conducting layer is disposed at least partially on the second surface and at least one opening in the essentially rigid substrate extends from the first surface to the second surface. An ink ejection assembly has a first surface and a second surface and includes an ejector. This ejector imparts momentum to ink within the ink ejection assembly. A nozzle plate forms a first portion of the first surface of the ink ejection assembly and has a nozzle through which ink is ejected. A conducting layer forms a second portion of the first surface of the ink ejection assembly so that the ejector may be energized. An ink feed channel is disposed in the second surface of the ink ejection assembly. The ink ejection assembly is disposed in the at least one opening and in contact with the essentially rigid substrate such that the nozzle plate and the first surface of the essentially rigid substrate are in a predetermined relationship relative to one another and the conducting layer of the ink ejection assembly is connected to the conducting layer of the essentially rigid substrate."}
{"text": "Digital signatures have been provided for non-repudiation wherein the user is associated with a given set of data and the signer can not deny participation with the signature. A digital signature associates a digital or numerical code with a set of electronic data. The code is generated using a private key that uniquely identifies the person that is approving the contents of the data.\nTo create a digital signature, an encryption process, such as DES or Triple DES for example, is utilized with a private key to encrypt a hash of the data set. The private key is associated with a given user. Thereafter, a public key which corresponds to the private key is utilized to decrypt the encrypted data to arrive at the original data set.\nDigital signatures provide security benefits of identification and authentication plus data integrity to arrive at non-repudiation. The identity of the signing person of a transaction is known and can be proven to a third party. The signature is linked to the user. The signature is also linked to the data being signed such that if the data is changed, the signature is invalidated. For non-repudiation, the signing party can not deny having signed the transaction inasmuch as the signature is linked to the user and the data.\nPrevious paper-oriented tasks such as vendor payment were processed by routing paper around the organization. Approvals were hand-written signatures on paper. However, the aims of arriving at substantially paperless initiatives and improving the practice of obtaining authorization have resulted in wider acceptance and usage of digital signatures.\nThere exists a need to provide improved storage and verification systems utilizing digital signatures."}
{"text": "The present invention relates to a liquid crystal device for use in, e.g., a display apparatus for displaying images including characters and/or figures, particularly a liquid crystal device using a chiral smectic liquid crystal suitable for full-color display and a liquid crystal device having a stripe electrode structure suitable for a simple matrix driving. The present invention also relates to a process for producing the liquid crystal device and a color liquid crystal display apparatus using the liquid crystal device.\nA display device of the type which controls transmission of light in combination with a polarizing device by utilizing the refractive index anisotropy of ferroelectric (or chiral smectic) liquid crystal molecules has been proposed by Clark and Lagerwall (U.S. Pat. No. 4,367,924, etc.). The ferroelectric liquid crystal has generally chiral smectic C phase (SmC*) or H phase (SmH*) of a non-helical structure in a certain temperature region and, in the SmC* or SmH* phase, shows a property of assuming either one of a first optically stable state and a second optically stable state (bright and dark states) responding to an electrical field applied thereto and maintaining such a state in the absence of an electrical field, namely bistability, and also has a quick responsiveness to the change in electrical field. Thus, it is expected to be utilized in a high speed and memory type display device and particularly to provide a large-area, high-resolution display based on its excellent function.\nFIG. 1 shows a sectional view of a liquid crystal device using a ferroelectric liquid crystal based on two-valued (white and black) display.\nReferring to FIG. 1, the liquid crystal device (panel) includes insulating substrates 1a and 1b, transparent electrodes 6a and 6b, auxiliary electrodes 7a and 7b, short-circuit prevention layers 8a and 8b, roughened surface-forming layers 9a and 9b, alignment layers 10a and 10b, an adhesive bead 11 (after adhesion), a spacer bead 12, and a liquid crystal layer 13. Each of the transparent electrodes 6a and 6b constitutes drive electrodes in combination with the auxiliary electrodes 7a and 7b, respectively. The drive electrodes (including the electrodes 6a and 7a and the electrodes 6b an 7b, respectively) intersect with each other at right angles to form a matrix electrode structure. At each intersection, one pixel is constituted and corresponds to a region between two broken lines in FIG. 1.\nIn order to ensure a good impact (shock) resistance and keep a uniform thickness of a liquid crystal layer, an ordinary liquid crystal panel needs to use spacer beads 12 composed of an adhesive and softer material (than the spacer beads), i.e., for effecting adhesion between the upper and lower substrate members (structures) as described above.\nSuch spacer beads and adhesive beads have been dispersed together on one of the pair of substrate in an ordinary (liquid crystal) panel production process (substrate production process). However, in the subsequent steps, particularly in the step of applying the pair of substrates to each other, the adhesive beads are liable to be moved due to flowability and poor adhesiveness thereof at that time, thus adversely affecting performances of a resultant liquid crystal panel.\nFurther, the above-mentioned liquid crystal device using a ferroelectric liquid crystal has a very small cell gap (i.e., a thickness of a ferroelectric liquid crystal layer), so that the injection of the liquid crystal into the cell gap of a blank cell is not readily performed, thus resulting in a defective liquid crystal panel in a relatively high proportion. For this reason, the ferroelectric liquid crystal device is required to improve a production yield.\nThere has been known a liquid crystal device including the above ferroelectric liquid crystal device having a matrix electrodes structure such that a pair of substrates (electrode plates) each provided with a group of electrodes in the form of stripes are oppositely disposed so as to form a pixel at each intersection at right angles and a gap between the substrate is filled with a liquid crystal. In case where such a liquid crystal device causes a short-circuit between the electrodes and has an electrode resistance out of its specifications, it is almost difficult to repair or replace such defective electrodes. For this reason, in an actual production line, after the formation of upper and lower electrode groups, all of the electrodes are subjected to inspection (check) with respect to short-circuit and electrode (wire) resistance by placing an inspection terminal on each lead-out portion of the electrodes, thus removing defective products from the production line.\nAs the number of pixels per unit display area is increased for providing a higher definition display image, an electrode width for each stripe electrode becomes narrower (smaller). Accordingly, in this case, an inspection terminal is not readily placed on a lead-out portion of an objective electrode in the above-described inspection stage, thus being liable to fail to perform a correct inspection operation.\nFurther, in order to reduce the electrode (wire) resistance, a metal wire (metal layer) as an auxiliary electrode is generally formed on a transparent electrode within an extent not impairing a display quality. The metal wire is liable to be damaged (e.g., burned out) by an inspection terminal having a narrowed top portion corresponding to a small electrode width or is liable to cause short-circuiting with a metal piece (fragment) scraped off or removed by the terminal.\nIn case where the liquid crystal device as described above is incorporated into a color liquid crystal display apparatus, a color filter comprising color filter segments of at least three colors including red (R), green (G), blue (B), and optional transparent color (W: white) in the form of stripes or a mosaic color filter wherein any adjacent (parallel) two color filter elements (comprising R, G, B and optional W segments) in one direction are shifted from each other by xc2xd pitch of one color filter segment in the direction may generally be used. Such a color filter is generally disposed at an inner surface of one of upper and lower (a pair of) glass substrates (i.e., on a side closer to a liquid crystal layer), whereby a resultant liquid crystal device has different layer structures with respect to the upper and lower substrates different from the case of a monochromatic (white and block) liquid crystal display apparatus.\nIn another aspect, a chiral smectic liquid crystal (e.g., a ferroelectric liquid crystal or an anti-ferroelectric liquid crystal) shows one orientation (alignment) state under application of an electric field of one polarity based on a certain reference potential level and shows the other orientation state under application of an electric field of the opposite polarity. Such a property is quite different from that of a twisted nematic (TN)-type liquid crystal. There has been developed a color liquid crystal display apparatus utilizing the above property of the chiral smectic liquid crystal.\nThe above-mentioned two orientation states of the chiral smectic liquid crystal are required to have potential energies having symmetry. However, if a pair of (upper and lower) substrates have different layer structures thereon from each other as described above, the potential energies of the two orientation states are liable to become asymmetrical. The asymmetry of the potential energies is liable to cause that of switching threshold values between orientation state and the other orientation state.\nThe above problem is peculiar to the chiral smectic liquid crystal and does not substantially arise in the case of the TN liquid crystal. Particularly, the asymmetry of switching threshold values is liable to narrow (decrease) a drive margin (a margin allowing a good display state) determining a latitude in selecting drive signal waveform conditions, such as a voltage level, a pulse width and a frequency.\nAccordingly, in the color liquid crystal display apparatus, it is important to find out a parameter (or factor) largely affecting the drive margin and to appropriately select and control such a parameter for providing a wider drive margin.\nA first object of the present invention is to provide a liquid crystal (particularly a chiral smectic liquid crystal) device allowing a uniform liquid crystal layer thickness and a good shock (impact) resistance to retain good panel performance and improved in production yield and capable of realizing a full color display apparatus having high qualities comparable to those of a display apparatus using a cathode ray tube (CRT).\nA second object of the present invention is to provide a liquid crystal device having a electrode structure capable of readily ensuring inspection (check) regarding an occurrence of short-circuit and an electrode (or wire) resistance without damaging electrodes used.\nA third object of the present invention is to provide a color liquid crystal display apparatus capable of effecting good display in any operation conditions while retaining a wider drive margin.\nAccording to the present invention, the above first object is principally accomplished by a liquid crystal device, comprising: a pair of substrates each provided with an electrode including one substrate having thereon a color filter and a coating layer, and a liquid crystal layer comprising a chiral smectic liquid crystal disposed together with spacer beads between the pair of substrates, wherein\nthe liquid crystal layer has a thickness smaller than a diameter of the spacer beads and a maximum thickness of the coating layer, the coating layer having a pencil hardness of at most 7 H.\nThe above objects of the present invention are accomplished by a process for producing a liquid crystal device, comprising the steps of:\nforming on a first insulating substrate a light-interrupting layer, a color filter comprising plural color filter segments, a coating layer, a barrier layer, a transparent electrode, an auxiliary electrode, a short-circuit prevention layer, a roughened surface-forming layer, and an insulating layer in succession in this order,\nforming on a second insulating substrate a transparent electrode, an auxiliary electrode, a short-circuit prevention layer, a roughened surface-forming layer, and an insulating layer in succession in this order,\nrubbing the surface of each of the insulating layers on the first and second substrates, dispersing adhesive beads over the alignment layer surface formed on the first substrate or the second substrate,\ndisposing a sealing agent having a prescribed pattern on the insulating layer surface formed on the second substrate or the first substrate,\ndispersing spacer beads over the alignment layer surface provided with the sealing agent,\nadhesively bonding the first and second substrate to each other while fixing the first or second substrate over which the adhesive beads are dispersed,\nscribing the first and second substrates to remove an unnecessary portion,\ninjecting a chiral smectic liquid crystal from an injection port into a gap between the first and second substrates, and\nsealing up the injection port.\nAccording to the present invention, the above second object is principally attained by a liquid crystal device, comprising: a pair of oppositely disposed substrates each provided with a group of transparent electrodes in the form of stripes, and a liquid crystal disposed between the substrates, wherein\neach of the transparent electrodes partially has an auxiliary electrode in its length direction and has both lead-out end portions in a region other than a display region, each of the lead-out end portions including an exposed check portion where the auxiliary electrode is patternized so as to expose the transparent electrode.\nIn this case, the exposed check portions may preferably have a width larger than that of the remaining portion and are disposed alternately at every transparent electrode in their width direction.\nFurther, each of the exposed check portion may preferably have both end portions where the auxiliary electrode is connected so as to enclose the exposed check portion.\nThe device may preferably include a dummy electrode in a region other than a display region so as to form a pattern similar to those of groups of data electrodes and scanning (common) electrodes, whereby measurement of an electrode resistance is performed without adversely affecting the data and scanning electrodes.\nAccording to the present invention, the above third object is principally achieved by a color liquid crystal display apparatus, including:\na liquid crystal device, comprising: a pair of oppositely disposed first and second substrates each provided with a group of transparent electrodes in the form of stripes, and a liquid crystal layer comprising a chiral smectic liquid crystal disposed together with spacer beads between the pair of substrates, the first substrate having thereon a color filter comprising plural color filter segments and a coating layer, wherein the transparent electrodes on the second substrate have a width smaller than that of the transparent electrodes on the first substrate,\nscanning signal supply means for supplying scanning signals to the transparent electrodes on the first substrate, and\ndata signal supply means for supplying data signals including an interval at a prescribed temperature or below to the transparent electrodes on the second substrate, each of the data signals corresponding to each of color filter segments of the color filter.\nThe above color liquid crystal apparatus is effective in ensuring a wide drive margin and providing a good display state in any environmental conditions by appropriately setting and fixing driving conditions allowing the wide drive margin based on its structural characteristic features.\nThe present invention further provides a liquid crystal device, comprising: a pair of oppositely disposed substrates each provided with a group of transparent electrodes, and a liquid crystal disposed between the substrates, the groups of transparent electrodes of the pair of substrates intersect with each other to form a pixel at each intersection, wherein\none of the substrates includes a light-interrupting layer for covering a part of the pixel located in a position corresponding to at least one end portion of the pixel.\nThese and other objects, features and advantages of the present invention will become more apparent upon a consideration of the following description of the preferred embodiments of the present invention taken in conjunction with the accompanying drawings."}
{"text": "Today's hospitality systems usually support a subscription of video or audio information like movies, music etc. which are displayed by an in-room screen or played by in-room loudspeakers.\nA disadvantage of current prior art hospitality systems is the limited way to deal with preferences of individual guests regarding the lighting of the guest's room, temperature etc.\nYet another disadvantage of known hospitality systems is that an identical selection of information and entertainment contents is offered to all guests. Every guest is encountered by a bulk of selectable entertainment media and other contents which will limit the usability of selecting media in the guest's sphere of interest.\nA further disadvantage of known hospitality systems is the limited manner to report alarm situation to guests and to effectively guide personnel and guests in the case of an evacuation."}
{"text": "Sachets of flexible film forming material are known. In the food packaging industry it has been conventional to \"form, fill and seal\". That is the individual flexible container is formed, it is then filled with food product and the container sealed. This means that one filling head can fill only one container at a time, For small containers production times are increased because the time taken to bring the container into alignment with the filling head and then to remove it, is significant and the volume of product filled from any filling head is correspondingly reduced when smaller containers are filled.\nFor consumer use and for applications where measured quantities of product are desirable in disposable containers there is a demand for small volume containers.\nThe cost to food packagers of using a form, fill and seal method as opposed to buying preformed packages and only filling and sealing is also significant. However preformed packages especially for small volumes are difficult to handle. Further where aseptic techniques are used sterilization becomes a problem. Thus in these applications the form, fill and seal method has been seen to be advantageous because a sterilization step can be carried out between the forming and filling sequences.\nA number of proposals have been made to increase the volume of production of small containers using the form, fill, seal method. U.S. Pat. No. 4,517,787 is a recent example and U.K. No. 2054511 is an example of forming and filling a plurality of small compartments which are subsequently separated to form individual packages."}
{"text": "Many drugs operate by chemically binding to specific molecular receptors. Molecular receptors typically are specific proteins (a term that includes glycoproteins and lipoproteins) in an animal such as a human being, and drug design and selection can be facilitated by accurately estimating the binding affinity of a drug to a protein, or, more generally, estimating the binding affinity of a ligand to a receptor, the term receptor being used to mean any moiety that specifically binds the ligand.\nOne way to determine receptor-ligand binding affinity uses the molecular structure of the receptor-ligand complex that results when the ligand binds to the receptor. Such structures may be studied by x-ray crystallography. The publicly accessible protein data bank (PDB) now contains more than 10,000 x-ray crystal structures, and pharmaceutical and biotechnology companies have an order of magnitude more proprietary structures. Many of these structures have been co-crystallized with small molecules bound to them. The examination of such structures, and deployment of the knowledge thereby gained to design new, more potent, and more specific inhibitors, is referred to as structure-based drug design.\nComputational modeling facilitates structure-based drug design. One aspect of modeling detailed below involves scoring functions that use simulation techniques, such as molecular dynamics, Monte Carlo, or continuum electrostatics calculations. These scoring functions can be problematic, as one is required to calculate a very small difference (the binding affinity) between two very large numbers (the free energies of the complex and of the separated protein and ligand). An alternative approach is to develop an empirical scoring function, based on the geometry of the complex, which directly evaluates the desired quantity. Such an approach has the advantage of being extremely fast as well as being amenable to fitting to experimental data, large amounts of which are now publicly available. Our own application, U.S. patent application Ser. No. 11/373,684 (US 2007/061118) entitled “Predictive Scoring Function for Estimating Binding Affinity,” which is hereby incorporated by reference in its entirety, discloses scoring functions.\nIt is desirable to increase the accuracy and robustness of scoring functions. It is particularly desirable to use scoring functions not only to distinguish active from inactive ligands, but also to rank ligands in a group of ligands, e.g., from most active to least active, to identify more promising candidates to be studied ahead of less promising candidates. Moreover, the ability to rank order a group of ligands having different chemotypes would allow one to search a broader range of candidates for effective drugs. Finally, rank ordering ligands without a large amount of preliminary work (e.g., using a training set) would further increase the convenience and speed of the process."}
{"text": "For over four decades, aminoacyl-tRNA synthetases (AARSs) were thought of as essential housekeeping proteins that catalyze the aminoacylation of tRNA molecules as part of the decoding of genetic information during the process of protein translation. AARSs have been extensively studied in this respect, and many of their full-length sequences were cloned for sequence analysis and to provide a rich source of biochemical experimentation. Some fragments of AARSs, and other proteins, however, possess unexpected activities not associated with aminoacylation, including extracellular signaling activities that modulate pathways beyond protein translation. Generally, these unexpected activities are not observed in the context of the full-length or parental protein sequences; instead, they are observed following removal or resection of AARS protein fragments from their parental sequences, or by expressing and sufficiently purifying fragment AARS sequences and then testing for novel, non-synthetase related activities.\nWhile the full-length sequences of AARS have been known for some time, no systematic experimental analysis has been conducted to elucidate such AARS protein fragments, or protein fragments from related or associated proteins, or to evaluate the potential role of the full length AARS proteins for novel biological activities outside of the context of amino acid synthesis. In portions of this specification, such AARS protein fragments, AARS domains, or AARS alternative splice variants are referred to herein as “resectins”. In its broadest context, the term “resectin” refers to a portion of a protein which has been excised or restricted (either by means of proteolysis, alternative splicing, mutagenesis, or recombinant genetic engineering) from the context of its native full-length or parental protein sequence, which often otherwise masks its novel biological activities. Likewise, no systematic experimental analysis has been conducted to explore the use of such resectins as biotherapeutic agents, diagnostic agents, or drug targets in the treatment of various medical conditions, or their potential association with human diseases. As essential housekeeping genes with a known function in mammals that is critical to life, AARSs were neither considered as drug targets in mammals, nor were they parsed out by standard genomic sequencing, bioinformatic, or similar efforts to identify resectins having non-synthetase activities. Standard biochemical research efforts have similarly been directed away from characterizing the biological properties of AARS resectins and their potential therapeutic and diagnostic relevance, mainly due to the previously understood role of their corresponding full-length parental AARSs."}
{"text": "1. Field of the Invention\nThe present invention pertains to liquid crystal on silicon (LCOS) displays, and more particularly to improved pixel cell design for liquid crystal on silicon displays with enhanced voltage control and simplified circuit to achieve prior to display frame data loading.\n2. Description of the Prior Art\nLiquid crystal on silicon (LCOS) microdisplay technology is still challenged by the need to DC balance the liquid crystal material accurately while generating images that are free of flicker and while limiting the RMS voltages to the useful range of the electro-optic efficiency curve of the liquid crystal device. Most liquid crystal devices disclosed in the art fail to meet one or more of these challenges. The current LCOS devices may be thought of as being divided into three classes, based on the method of creating grayscale. Each is prone to a particular class of problems. The three classes of devices are analog drive, digital drive based on simple SRAM, and complex digital drive. All are prone to some degree to the related problems of DC offset and image flicker. The degree to which these problems manifest themselves can influence product acceptability and product longevity and therefore solutions that mitigate or eliminate these problems are worthy subjects of invention.\nThere are additional considerations beyond the basic problems cited above. For example, in order to achieve savings on power consumption and prolong the life of a display system, it is desirable to have a way of inverting the voltage applied to the liquid crystal pixel without having to change the state of the memory cell and without having to rely on the state of the logic memory cell in order to directly supply a voltage onto the pixel electrode. However, the conventional multiplexing devices that attempt to narrow the range of voltages that are applied to the pixel electrode on an instantaneous basis and that receive input signals directly from a memory cell fail to provide the needed flexibility because they are not capable of independently controlling the memory state that gets driven through the multiplexer to the pixel mirror. Limited by these technical difficulties, the conventional technologies of LCOS display are provide displays of higher quality only with difficulty. Specifically, the displays are often hindered by problems of image sticking and flicker due to the low DC balancing rates as will be further explained below.\nLiquid crystal display (LCD) technology has progressed rapidly in recent years, and has become an increasingly common option for display systems. LCD's make up the largest portion of the flat panel display market. This market dominance is expected to continue into the future. The superior characteristics of liquid crystal displays with regard to weight, power, and geometry in image visualization, have enabled them to compete in fields historically dominated by Cathode Ray Tube (CRT) technology, such as high definition television systems, desktop computers, projection equipment, and large information boards. As the cost of LCD systems continues to fall, i.e., is predicted that they will eventually take over the market for traditional CRT applications.\nThe biggest disadvantages of current CRT systems are their bulky size, geometry, and weight, as well as their high power consumption. These disadvantages are clearly evident when comparing the features of CRT and LCD projection displays with similar characteristics. In general, projection display systems offer several additional advantages over CRT systems. First, projection display systems offer the possibility of using large screens for group viewing with the ability to easily change the image size and position. Second, projection display systems offer high performance, and the ability to accept image data input from a variety of devices such as computers, television broadcasts, and satellite systems. Virtually any type of video input can be projected through such a system. The application of LCDs to projection systems has further attractive features such as high brightness, high resolution, and easy maintenance. LCD front projection displays provide higher resolution and brightness than comparable CRT-based systems. In comparison with CRT's, installation of LCD projection systems is easy and their viewing angles are generally much wider. Most front projection LCD display systems are compatible with personal computers and can operate with video signals from television systems. LCD front projectors are easily adapted for applications such as home theaters.\nTypically, LCD projection systems include small LCD panels, usually ranging from 0.5 to 5 inches in diagonal, a series of dichroic mirrors or filters, and a series of projection lenses to cast the images onto a screen. Commonly, three panel systems are used, where one or more dichroic mirrors divide white light coming from a light source, into the three primary colors of red, green, and blue (RGB). The dichroic minors direct each of the RGB components toward a separate LCD panel. The corresponding LCD panel modulates each of the RGB components of the light according to the input image data corresponding to that color. Output dichroic minors synthesize the modulated RGB light components and project the image onto a viewing screen.\nTo enhance the luminance and fill factor of the liquid crystal projection panels, reflective LCD pixels are often used. These systems, referred to as Liquid Crystal on Silicon micro-displays (LCOS), utilize a large array of image pixels to achieve a high-resolution output of the input image. Each pixel of the display includes a liquid crystal layer sandwiched between a transparent electrode and a reflective pixel electrode. Typically, the transparent electrode is common to the entire display while the reflective pixel electrode is operative to an individual image pixel. A storage element, or other memory cell, is mounted beneath the pixels and can selectively direct a voltage on the pixel electrode. By controlling the voltage difference between the common transparent electrode and each of the reflective pixel electrodes, the optical characteristics of the liquid crystal can be controlled according to the image data being supplied. The storage element can be either an analog or a digital storage element. More and more often, digital storage elements, in the form of static memory are being used for this purpose.\nThe liquid crystal layer rotates the polarization of light that passes through it, the extent of the polarization rotation depending on the root-mean-square (RMS) voltage that is applied across the liquid crystal layer. (The incident light on a reflective liquid crystal display thus is of one polarization and the reflected light associated with “on state” is normally of the orthogonal polarization.) The reason that the degree of polarization change depends on the RMS voltage is well known to those skilled in the art—it is the foundation of all liquid crystal displays.\nTherefore, by applying varying voltages to the liquid crystal, the ability of the liquid crystal device to transmit light can be controlled. Since in a digital control application, the pixel drive voltage is either turned to dark state (off) or to light state (on), certain modulation schemes must be incorporated into the voltage control in order to achieve a desired gray scale that is between the totally on and totally off positions. It is well known that the liquid crystal will respond to the RMS voltage of the drive waveform in those instances where the liquid crystal response time is slower than the modulation waveform time. The use of pulse-width modulation (PWM) is a common way to drive these types of digital circuits. In one type of PWM, varying gray scale levels are represented by multi-bit words (i.e. a binary number) that are converted into a series of pulses. The time averaged RMS voltage corresponds to a specific voltage necessary to maintain a desired gray scale.\nVarious methods of pulse width modulation are known in the art. One such example is binary-weighted pulse-width-modulation, where the pulses are grouped to correspond to the bits of a binary gray scale value. The resolution of the gray scale can be improved by adding additional bits to the binary gray scale value. For example, if a four-bit word is used, the time in which a gray scale value is written to each pixel (frame time) is divided into fifteen intervals resulting in sixteen possible gray scale values (24 possible values). An 8-bit binary gray scale value would result in 255 intervals and 256 possible gray scale values (28 possible values).\nSince most nematic liquid crystal materials only respond to the magnitude of an applied voltage, and not to the polarity of a voltage, a positive or negative voltage, of the same magnitude, applied across the liquid crystal material will normally result in the same optical properties (polarization) of the liquid crystal. However, the inherent physical characteristics of liquid crystal materials cause deterioration in the performance of the liquid crystal material due to an ionic migration or “drift” when a DC voltage is applied to them. A DC current will cause the contaminants always present in liquid crystal materials to drift toward one alignment surface or the other, if the same voltage polarity is continuously applied. This will result in the contaminants plating out onto the alignment layer with the result in that the liquid crystal material will begin to “stick” at an orientation and not respond fully to the drive voltages. This effect is manifested by the appearance of a ghost image of the previous image that is objectionable to viewers. Even highly purified liquid crystal materials have a certain level of ionic impurities within their composition (e.g. a negatively charged sodium ion). In order to maintain the accuracy and operability of the liquid crystal display, this phenomenon must be controlled. In order to prevent this type of “drift”, the RMS voltage applied to the liquid crystal must be modified so that alternating voltage polarities are applied to the liquid crystal. In this situation, the frame time of the PWM is divided in half. During the first half of the frame the modulation data is applied on the pixel electrode according to the predetermined voltage control scheme. During the second half of the frame time, the complement of the modulation data is applied to the pixel electrode. When the common transparent electrode is maintained at its initial voltage state, typically high, this results in a net DC voltage component of zero volts. This technique generally referred to as “DC Balancing” technique is applied to avoid the deterioration of the liquid crystal without changing the RMS voltage being applied across the liquid crystal pixel and without changing the image that is displayed through the LCD panel.\nModulation schemes that are employed to drive the liquid crystal pixel elements must therefore be able to accurately control the amount of time the pixel on and “off”, in order to achieve a desired gray scale from the pixel. The degree of rotation of light that occurs follows the RMS voltage across the liquid crystal pixel. The degree of rotation in turn affects directly the intensity of the light that is visible to the observer. In this manner modulating voltages influences the intensity perceived by an observer. In this manner gray scale differences are created. The combination of all of the pixels in a display array results in an image being displayed through the LC device. In addition to controlling the root-mean square (RMS) voltage that applied to the pixel, the polarity of the voltage must be continuously “flipped-flopped” so that deterioration of the liquid crystal is avoided. Known modulation schemes are not able to prevent liquid crystal deterioration while still being able to accurately control the RMS voltage as that applied to the liquid crystal in many projection applications where the display device is subjected to high heat loading and high light intensity. This results from a combination of electron-hole-pair generation and increased electron mobility due to the aforementioned conditions.\nOther liquid crystal display system, more particularly those in flat panel TFT displays, drive the liquid crystal at slower frequencies on the order of 60 Hz. These systems create gray scale by controlling the charge placed on a storage capacitor that is connected to a transparent pixel electrode that forms part of the pixel element. Increasing or decreasing the charge on the storage capacitor is then applied to realize a gray scale effect. These devices typically refresh the storage capacitor at a rate on the order of 60 Hz to 90 Hz. In this case the liquid crystal may not be responding to the RMS voltage of the display but rather directly to the DC charge placed on the device. This is because the time constants of the liquid crystal are often actually faster than the refresh rate on the device. In this case the display is extremely sensitive to ionic contamination because the charge on the pixel electrode will bleed off quickly and the display will flicker if such ionic contamination is present. The common figure of merit that is used to describe ionic contamination for such displays is the voltage holding ratio or charge holding ratio, a high percent holding ratio indicating low contamination.\nThe electro-optical properties of many liquid crystal materials cause them to produce a maximum brightness at a certain RMS voltage (VSAT), and a minimum brightness at another RMS voltage (VTT). Applying an RMS voltage of VSAT results in a bright cell, or full light reflection, while applying an RMS voltage of VTT results in a dark cell, or minimal light output. Increasing the RMS voltage to a value above that of VSAT, may reduce the brightness of the cell rather than maintaining it at the full light reflection level. Likewise decreasing the RMS voltage to a value below that of VTT, will normally increase the brightness of the cell rather than maintaining it at the zero light reflection level. At RMS voltages between VSAT and VTT the percent brightness increases (or decreases, depending on the electro-optic mode) as the RMS voltage increases. The voltage range between VTT and VSAT therefore defines the useful range of the electro-optical curve for a particular liquid crystal material. It follows that RMS voltages outside of this range are not useful and will cause gray scale distortions if applied to the crystal pixels. It is therefore desirable to confine the RMS voltages applied to the pixels to this useful range between VSAT and VTT. Many known display systems drive the logic circuitry with voltages that are outside of the useful range of the liquid crystal, and applying these voltages directly onto the pixel electrode results in wasted power. For example, logic circuitry may operate at 0 and 5 volts or 0 and 3.3 volts. If the useful range of the liquid crystal material is inside of this range, more time and power must be expended to achieve RMS voltages that are within the useful range. In a system that has a useful VTT to VSAT range of 1.0 to 2.5 volts and that has logic circuitry that operates at 0 to 5 volts, in order to achieve an RMS voltage of 2.5 volts, the pixel must see an equal amount of the 0 volt state and the 5 volt state over a time frame in order to achieve an RMS voltage of 2.5 volts. It would be much more efficient if the logic circuitry operated at the VSAT and VTT levels, rather than at levels outside of the VSAT to VTT range. This would make the time averaging simpler and faster and less power would be required to drive the same systems.\nFor these reasons, it is desirable to confine the RMS voltages to the useful range of the electro-optical response curve of the liquid crystal material. Furthermore, since the technique of rewriting data to achieve DC balance is used in most of the conventional analog driven TFT panels and micro-displays, it is also desirable to implement a D-C balancing technique for alternating the voltage applied on the liquid crystal, without having to continuously write new data onto the storage element. One of the techniques of accomplishing this it by inverting the voltage applied to the liquid crystal pixel without having to change the state of the memory cell and without having to rely on the state of the logic memory cell in order to direct a voltage onto the pixel electrode.\nConventional multiplexing devices are employed in attempt to narrow the range of voltages that are applied to the pixel electrode on an instantaneous basis and that receive input signals directly from a memory cell. But since the entire panel is written with voltage scales that alternate between a higher voltage range and a lower voltage range, conventional multiplexing devices fail to provide the needed flexibility because they are not capable of independently controlling the memory state that gets driven through the multiplexer to the pixel mirror. For example, in “Miniature FLC/CMOS Color Sequential Display Systems”, SID Digest, 1997, Section 21.3, Handschy et al describe a backplane based on an SRAM device with full static control logic and row and column drivers. In the particular instance the pixel mirror voltage is determined solely by the logic state of the SRAM circuit underlying it. Thus the pixel may receive 0 volts or 5 volts, depending on its logic state. The common electrode value is set at 2.5 volts. No provision is described for decoupling DC balance from the writing of data and the device appears to be a simple SRAM device modified by the addition of pixel mirrors and post processing to make the device part of a liquid crystal cell. A further description, of this is found in PCT Publication WO 01/16928. Xue, et al, “Reduction of Effects Caused by Imbalanced Driving of Liquid Crystal Cells.” While simple, the device does not enable the enhanced performance that comprises part of the present invention.\nAnother example of display system is disclosed in U.S. Pat. No. 6,005,558. A display system includes a memory element coupled to a multiplexer. Depending on the state of the memory element, the multiplexer directs one of two predetermined voltages onto a pixel electrode. The multiplexer is situated externally to the memory cell and is controlled by external circuitry to operate in conjunction with DC balance and data load operations. In the disclosed invention, operation of the multiplexer external to the cell requires that the voltages delivered via the rails to the cell be modulated to provide DC balance. This adds substantially to the complexity of the device because the modulated voltage must be correct in all respects as these same voltages are used to drive the pixel mirrors and thus achieve DC balance. Design of a line that cam propagate a number of different voltages across long lines that must accurate in all cases is a significant design constraint. Furthermore, the disclosed invention requires that all elements be globally addressed to function. All these technical difficulties limit the effectiveness of the above inventions in providing practical solutions to the above-mentioned limitations.\nFurthermore, since the conventional systems utilize the state of the memory cell as a control signal to direct a voltage onto the pixel mirror, there is no independent means for selecting and directing a narrower range of voltages to be consistent with the electro-optical response curve, onto the pixel electrode. It often leads to further difficulties and limitations. Since the gray scale is symmetrical within those two voltage ranges on either side of the voltage of the common plain, there is a transition phase during data load where both upper and lower voltage ranges are present on the display at the same time. The duration of this period is perhaps 200 to 300 microseconds. Such approach places several limitations on the device: First, the common plane (VCOM) must be a fixed value because during the transition time data is present on the display in both upper and lower voltage ranges. This in turn means that the display must have a voltage authority range sufficient to permit the writing of analog data in both the upper and lower ranges. While occasionally this may be done within normal voltages, it is more often the case that the developer of the design must use special silicon processing techniques such as those associated with flash memories or with EEPROM technology. In the cases when the EEPROM is used, the resulting parts may have a 10–15 volt operating range so that Vcom can be set at the 4.5 to 6.5 volt range. This in turn provides a reasonable range of authority for symmetrical driving of the liquid crystal. In the former case, a normal silicon processes may be applied, the best CMOS processes can reach ˜6 volts, which dictates a VCOM of approximately 2.9 volts. However, with this voltage range, there are very few liquid crystal materials available to satisfy the requirement of the panel. The best design, manufacturing and yield economics are associated with standard CMOS designs rather than the high voltage processes. For these reasons, manufacturing LCOS display panels and other LCD devices often results in a lower yields than average yield rate. Thus, the production costs are increased when non-standard CMOS processes have to be carried out for manufacturing the LCOS device in order to satisfy these operational conditions when suitable liquid crystal materials are used.\nFurthermore the technique of applying a DC balance switching rate of once per data load in the conventional micro-displays further creates a situation where the ion migrating within the liquid crystal material begin to plate out toward the end of the data frame at a lower DC balancing rate. Eventually such displays begin to show more “image sticking.” where a ghost image reflecting older data remains in a display after new (and different) data is written to the display. It is objectionable and in fact is a specification item for TFT displays that the old data must dissipate within a specified short period. With a low DC balancing rate, a LCOS display often generates a second objectionable flicker artifact due to slow data rates caused by the existence of DC offset mechanisms. The manifestation is that VCOM can no longer be set to a point half way between the two voltage drive ranges but rather must be raise or lowered from that point to achieve a flicker-free image. This second problem creates a dilemma between solving flicker and eliminating image sticking. Further details are disclosed in a patented disclosure of U.S. Pat. No. 6,424,330 to Johnson, entitled “Electro-optic Display Device with DC Offset Correction” and the disclosures in that Patent is hereby incorporated by reference in this Application.\nA co-pending patent application Ser. No. 10/329,645 filed by an inventor of this Application discloses a pixel display configuration by providing a voltage controller in each pixel control circuit for controlling the voltage inputted to the pixel electrodes. The controller includes a function of multiplexing the voltage input to the pixel electrodes and also a bit buffering and decoupling function to decouple and flexible change the input voltage level to the pixel electrodes. The rate of DC balancing can be increased to one KHz and higher to mitigate the possibility of DC offset effects and the image sticking problems caused by slow DC balancing rates. The co-pending Application further discloses an enabling technology for switching from one DC balance state to another without rewriting the data onto the panels. Therefore, the difficulties of applying a high voltage CMOS designs are resolved. Standard CMOS technologies can be applied to manufacture the storage and control panel for the LCOS displays with lower production cost and higher yields. The DC-balancing controller in the co-pending Application is implemented with a ten-transistor (10-T) configuration consisted of two P-MOS transistors. While the controller can be more conveniently implemented, it does have a technical limitation due to a constraint that the P-MOS transistors are not effective in pulling down the voltage of the pixel mirror. The lower voltage limit V0 of the controller is set to 1.0 to 1.3 volts above the semiconductor ground voltage VSS depending on the design details of the circuits. The limitation occurs due to the fact that a P-MOS transistor is strong in pulling the voltage up to Vdd while weak in pulling down the voltage to VSS. The technical limitation of not able to operate the pixel control voltage closer to VSS thus reduces the flexibility of material selections of the liquid crystal materials for the LCOS image displays. In many cases, this more limited operation voltage range further reduces the response time of the liquid crystal. The image contrast is also adversely affected due to the reduced range of voltages for dark stage selection.\nIn addition to the above limitations, for a microdisplay system, in order to achieve a global or time-synchronously update of the image, it is required to write a frame of data to the buffer before the beginning of the display of that frame of data. The timing and sequences of loading the display data into the buffer would allow the image timing management to work as a set of frames rather than with a time skew between rows. Usually, it is optimal to allow the loading of the entire frame of data plus some small margin before the display of that frame to take into account of the minor timing discrepancies. For the purpose of achieving these goals, several conventional circuit designs employ a second stage device that can be implemented as capacitor subject to bleed off of charge during the display time or some sorts of secondary storage device that includes many transistors, e.g. 6T storage device. The secondary stage circuits add to the complexity and density of the display system and increase the probability of failure in the CMOS manufacturing processes. The added circuits further increases the size and volume of the display systems thus requiring a larger minimum circuit size and also increase the power consumptions and further increase the burdens of heat dissipation. Conventional technologies do not provide a effective solution to this limitation.\nFor these reasons, there is still need in the art of LCOS display to provide improved system with simplified circuit configuration without requiring the second stage circuit to achieve the loading to the data into the buffer before the display of the data. It is further desirable that the circuit can be driven either in a global update mode in a regular load mode that that the circuit designs can be implemented in more than one types of microdisplay systems."}
{"text": "This invention relates to process variable transmitters receiving power over two of three wires and communicating over a third wire to a controller."}
{"text": "Recently the incidence of serious fungal infection has become very prominent in patients undergoing chemotherapy for cancer, organ transplants and patients with AIDS. Most of these infections are caused by opportunistic pathogens like Candida spp., Aspergillus spp., Pneumocystis camii and Cryptococcus neoformans. The antifungal agents available in the market suffer with draw backs such as, toxicity, narrow spectrum of activity, fungistatic profile rather fungicidal. Some of them also exhibit drug-drug interactions and, as a result, therapy becomes very complex. In view of the high incidence of fungal infections in immunocompromised patients and the recent trend for the steady increase of the populations of these patients, demands for new antifungal agents with broad spectrum of activity and good pharmacokinetic properties have increased.\nWithin the available drugs to treat fungal infections, the azole class appears to be more promising. This class of compounds inhibit the biosynthesis of ergosterol in fungi, which is the main constituent of fungal cell membrane. Fluconazole and itraconazole are routinely used for maintenance of fungal infections. Although fluconazole is highly bioavailable, it is not active against filamentous fungi and emergence of fungal resistance has been reported recently (Antimicrob. Agents Chemother. 1995, 39, 1-8). Itraconazole is active against filamentous fungi, but it shows inconsistent results, maybe due to its high protein binding properties and less bioavailability. During the last few years, several research groups have been actively searching for new azoles with optimum pharmacokinetic properties. As a result, a number of candidate azoles have emerged, and some of them are undergoing preclinical and clinical evaluation. Some of the candidate azoles are disclosed in the following publications:\nSch 51048 (Drugs of the Future, 1995, 20, 241-247). PA0 Sch 56592; Antimicrob. Agents Chemother. (1996, 40, 1910-1913; 36th Interscience Conference Antimicrob. Agents Chemother. September 1996, New Orleans, Abst. F87-F102). PA0 UK-109, 496 (Drugs of the Future, 1996, 21, 266-271; EP 440372). PA0 TAK-187; 36th Interscience Conference Antimicrob. Agents Chemother. September 1996, New Orleans, Abst. F74; EP 567982). PA0 KP-103 (36th Interscience Conference Antimicrob. Agents Chemother. September 1996, New Orleans, Abst. F78, WO 94126734). PA0 ER-30346 (Drugs of the Future, 1996, 21, 20-24)\nIn the present invention, we report new triazoles with broad spectrum anti-fungal activity. The triazoles are particularly effective against systemic and lung invasive fungal infections."}
{"text": "Many buildings include complex systems to provide heating, cooling, lighting, security and other services. These systems may require a significant amount of energy, which can present a significant challenge to facility managers facing rising (and volatile) energy prices, as well as requirements to be more environmentally sensitive. These challenges can be particularly acute for large commercial buildings where manual control over many disparate systems is often time-consuming, burdensome, and expensive.\nThese and other issues are addressed by embodiments of the present disclosure. Among other things, embodiments of the disclosure can reduce energy costs and help control and manage building systems more efficiently and effectively compared to conventional systems or methods for building control."}
{"text": "Water exposure is among major reasons that may cause significant malfunction of devices, such as electronic devices, which include electronic, mechanical, or chemical components. Therefore, verification of significant water exposure (or water immersion) is important to manufacturers of the devices. For example, for purposes such as warranty claim assessment, trouble-shooting for repairs, and product development, a manufacturer of electronic devices typically needs to verify whether a malfunctioning device has been immersed in water. FIGS. 1A-B illustrate an example conventional arrangement for verifying water immersion of a device 100 such as a cellular phone.\nAs illustrated in FIGS. 1A-B, device 100 includes a compartment 105 with a cover 108 for housing a removable module 110 such as a battery. For verifying water immersion, device 100 may include a detector 102 disposed at a bottom of compartment 105. To determine whether device 100 has previously been immersed in water, an inspector 170, such as a representative of the manufacture of device 100, may open cover 108 and remove removable module 110 to see whether detector 102 has changed color.\nDetector 102 may be commercially available from suppliers such as 3M Company (www.3m.com) of St. Paul, Minn. and Schreiner Label Tech, Inc. of Southfield, Mich., and NovaVision Inc. (www.novavisioninc.com) of Bowling Green, Ohio and typically has a configuration as illustrated in FIG. 1B.\nFIG. 1B illustrates a cross-sectional view of detector 102, with dimensions of detector 102 exaggerated for clear illustration. As shown in FIG. 1B, detector 102 includes an adhesive layer 132 for securing detector 102 at the bottom of compartment 105. Over adhesive layer 132, detector 102 includes a dye layer 112 that may dissolve after receiving a certain amount of water. Over dye layer 112, detector 102 includes a paper layer 122 for bearing dye layer 112 as well as allowing dissolved dye layer 112 to diffuse (or permeate) through paper layer 122 to show an indication (such as a red spot) that is visible to inspector 170 after detector 102 has been significantly exposed to (or immersed in) water. Over paper layer 122, detector 102 includes a transparent cover layer 142 for limiting exposure of paper layer 122 to humidity, thereby further avoiding false positives.\nNevertheless, the above-described conventional arrangement for verifying water immersion with detector 102 has disadvantages. The disadvantages may pertain to limited application, tampering of detection 102, and complicated inspection process, described as follows.\nFirstly, application of the conventional arrangement is limited to devices that are designed to be easily opened. However, subject to requirements such as industrial standards, design considerations, or manufacturing considerations, some devices cannot be opened without being dismantled or damaged. Therefore, water immersion of some devices cannot be verified using the conventional arrangement.\nSecondly, since detector 102 is readily visible and accessible with easy removal of removable module 110, detector 102 may be easily tampered. For example, detector 102 may be replaced with a new detector after device 100 has been immersed in water. As a result, it may become difficult for the manufacturer of device 100 to verify or prove water immersion of device 100.\nThirdly, the process of water immersion verification may be complicated. When the representative of the manufacturer is to inspect detector 120, the representative must remove cover 108 and removable module 110 before inspection. The representative may also need to make sure that cover 108 and removable module 110 are replaced back to device 100 after the inspection. Accordingly, the inspection process is complicated, and the inspection process may be further complicated when the representative handles multiple devices."}
{"text": "The present disclosure relates to bands for wearing around an ankle or a wrist. The present disclosure also relates to stands for portable electronic devices that may support a portable electronic device in a desired orientation for viewing, that may provide for wireless charging of a portable electronic device and for exchanging information with a portable electronic device. The present disclosure therefore relates to a multi-purpose charging interface for a portable electronic device, the charging interface being wearable on a user as a wristwatch or bracelet or anklet.\nElectronic devices, including portable electronic devices, have gained widespread use and may provide a variety of functions including, for example, telephonic, electronic messaging and other personal information manager (PIM) application functions. Portable electronic devices include, for example, several types of mobile stations such as simple cellular telephones, smart-phones, wireless personal digital assistants (PDAs), tablet computers, and laptop computers with wireless 802.11 or Bluetooth capabilities.\nIn today's world, a plethora of auxiliary devices have been developed to enhance a user's experience in using portable electronic devices. For example, a user may use a portable charging device for recharging a portable electronic device or may use an external speaker in connection with their portable electronic device. The user may employ an external interface for transferring information to or from the portable electronic device and may also employ an external display for viewing information transmitted from the portable electronic device. A stand or docking station may provide a convenient receptacle for an electronic device while coupling to a power source for charging a battery of the portable electronic device. Such stands may also provide an interface for data transfer between the portable electronic devices and, for example, a computer.\nUnfortunately, it can be inconvenient to reliably transport or locate such auxiliary devices when or where their use is desired."}
{"text": "1. Field of the Invention\nThe present invention relates to a maintenance apparatus for an inkjet-type print head.\n2. Description of the Related Art\nGenerally speaking, inkjet printers print by discharging ink drops to a desired position from plural nozzles in a print head mounted on a bidirectionally travelling carriage.\nMaintenance operations for appropriately cleaning the print head in a service area outside of the normal printing area include capping the print head with a cap that enables communication with air when the print head is not used for a long time, an ink vacuuming process for sucking ink that has increased in viscosity (referred to below as xe2x80x9cviscous inkxe2x80x9d) from inside the nozzles while the print head is sealed with the cap, and a wiping process using a wiper to wipe normal ink, viscous ink, and other contamination from the nozzle surface of the print head.\nDevices performing such maintenance operations must be able to move the cap and wiper toward and away from the print head. A maintenance apparatus according to the related art is therefore typically configured to move the cap in conjunction with movement of the carriage, and to move the wiper along a specific path through a cam mechanism using, for example, the drive power of a pump used for vacuuming ink as the drive power source, or is configured to move the cap and the wiper along separate paths using a similar cam mechanism.\nA problem with a maintenance apparatus according to the related art as noted above is that a longer carriage path must be provided in order for the cap to move in conjunction with carriage movement. This necessarily increases the size of the printer.\nThe cap must also be movable between three distinct positions: a retracted position where the cap is separated from the print head, a capping position where the cap covers the nozzle surface of the print head, and a sealed position where the nozzle surface is completely sealed for vacuuming ink from the nozzles. The cam mechanism required for the cap to move between these three positions independently of wiper movement is, however, complex and independent movement can be difficult to achieve.\nMore particularly, the space inside the cap must be able to communicate with the air while the cap covers the nozzle surface when in the capping position. A valve must therefore be provided in the cap, and a further problem is that the mechanism for opening and closing this valve is complex.\nThe present invention is directed to a solution to these problems, and an object of the invention is to provide a maintenance apparatus simplifying the cam mechanism for moving the cap and wiper and enabling the printer itself to therefore be made smaller.\nTo achieve these objects a maintenance apparatus for maintaining a print head having nozzles for discharging ink droplets and nozzle surface on which the nozzles are disposed, according to the present invention has a cap for covering the nozzles; a wiper for wiping the nozzle surface; and a cam member that is a rotatably disposed solid of revolution having on a side part thereof a first cam part for moving the cap and a second cam part for moving the wiper.\nBy linking movement of the cap and wiper using the first and second cam parts of the cam member, the present invention can cover or seal the print head with the cap without linking the capping mechanism to print head movement as in the related art. A printer comprising a maintenance apparatus according to the present invention can therefore be made smaller and the cam mechanism can be simplified because a single cam member is sufficient and a complicated cam mechanism is not required.\nFurther preferably, the maintenance apparatus also has a first slider movably supporting the cap and engaging the first cam part of the cam member to move the cap toward or away from the nozzle surface. A spring is also preferably positioned between the first slider and cap for urging the cap toward plane of the nozzle surface of the print head. The cap also preferably has a through-hole for communicating with air, and the first slider has a valve for opening and closing the through-hole. Yet further preferably, in this case, the first cam part of the cam member has a cam face for moving the cap between a sealed position whereat the valve is closed and the nozzles are covered, a covered position whereat the valve is open and the nozzles are covered, and a retracted position separated from the print head, in conjunction with cam member rotation.\nDriving a pump communicating with the cap when the cap is in the sealed position can purge ink inside the nozzles of the print head. When the cap is in the covered position the cap communicates with the air through the through-hole. Driving the pump in the covered position without vacuuming ink from the nozzles can therefore purge ink inside the cap. The first slider when pressed against the print head movably supports the cap, and the sealed and covered positions of the cap can therefore be set within the range of first slider movement. The cap can therefore be moved between the sealed position and covered position by simply changing the position of the slider, that is, by rotating the cam member.\nThe through-hole is preferably formed in the back of the cap and the valve is formed on the first slider at a position opposing the through-hole. The through-hole separates from the valve and the valve opens due to action of the spring when the cap moves from the sealed position to the covered position. The cap can therefore be easily changed from the sealed position to the covered position without using a complex valve mechanism.\nYet further preferably, the cam face of the first cam part has areas where the cap remains in each of the sealed position, the covered position, and the retracted position as the cam member rotates through a respective specific angle. This makes it easier to control movement of the cap to each of these positions.\nYet further preferably, the first slider has a lock part for fixing the print head in its home position. By making this lock part an integral part of the first slider the print head can be fixed in conjunction with movement of the first slider. It is therefore not necessary to provide a separate member for fixing the print head position and a mechanism for moving this separate member, and the capping mechanism itself is therefore simplified.\nYet further preferably, the second cam part of the cam member has a cam face for moving the wiper in conjunction with rotation of the cam member between a retracted position separated from the print head and plural wiping positions at different distances from the retracted position. By thus using a mechanism for changing the wiper position the wiper can be moved between, for example, a first wiping position for cleaning the nozzles and a second wiping position for cleaning the wiper itself, thereby enabling more precise maintenance. The cam face of the second cam part further preferably has areas where the wiper remains in each of the first wiping position, the second wiping position, and the retracted position as the cam member rotates through a respective angle. This makes it easier to control the movement of the wiper to each of these positions.\nThe first and second cam parts of the cam member are preferably related such that the wiper is in the retracted position when the cap is in the sealed position or covered position, and the cap is in the retracted position when the wiper is in a wiping position. The cap thus does not move during the wiping process, and the wiping operation can therefore be run independently of the capping process and ink vacuuming process.\nFurthermore, by making the maintenance apparatus smaller, the present invention also enables reducing the size of the printer, and the control components can also be simplified because the maintenance process can be accomplished by controlling primarily the angle of cam member rotation.\nOther objects and attainments together with a fuller understanding of the invention will become apparent and appreciated by referring to the following description and claims taken in conjunction with the accompanying drawings."}
{"text": "This invention relates to stator windings for miniature brushless dc motors in general and more particularly to an improved stator winding of the type in which several component windings are designed as injection molded frame windings.\nStator windings for a miniature brushless dc motor comprising a plurality of component windings designed as frame windings injection molded in plastic and nested together to form a cylinder surrounding the permanent magnet rotor of the motor and which can be inserted into a ring-shaped magnetic return yoke without slots are disclosed in German Offenlegungsschrift No. 1,613,379. In the stator winding disclosed therein the reaction moment acting on the winding is transmitted through a frictional connection between the cylinder and the magnetic return yoke which is connected to the housing. In order to achieve the necessary positive transmission between the return yoke and the cylinder formed by the frame windings, the cylinder must rest firmly against the return yoke. This requires a high dimensional accuracy in the cylinder which is to be inserted into the return yoke. However, when using parts which are injection molded in plastic such dimensional accuracy can be achieved only with great difficulty. A further problem arises because the plastic materials typically used in such applications change their dimensions as a function of climatic influences. The plastics can shrink or swell depending on the ambient temperature and air humidity. Because of these properties it is difficult to insert the cylinder formed by the frame windings into the return yoke in such a manner that perfect transmission of the reaction moment acting on the winding is assured at all times.\nThus, the need for an improved design for stator windings of this nature which will ensure reliable transmission of the reaction momemnt acting on the windings to the motor housing without placing any special requirements as to dimensional accuracy of the frame windings becomes evident."}
{"text": "1. Field of the Invention\nThis invention relates to an improved method for preparation of certain zeolites having a high silica/alumina mole ratio, i.e. greater than about 10. The improvement resides in adding a source of aluminum ions to a silica-rich amorphous reaction medium at a slow and controlled rate whereby the concentration of aluminum ions in the amorphous phase of the reaction mixture is maintained at steady state during crystallization.\nMore particularly, this invention relates to the above improved method of preparation of certain zeolites whereby crystallization time is substantially reduced from that required when conventional prior art methods of preparation are utilized and the resulting zeolite exhibits somewhat improved steam stability, i.e. stability toward steam deactivation and dealuminization.\nEven more particularly, this invention relates to organic compound conversion, such as, for example, catalytic conversion of oxygenates, such as methanol, and syn-gas conversion where water is always present during reaction, with the improved zeolite product of the present improved method as a catalyst.\n2. Discussion of Prior Art\nZeolitic materials, both natural and synthetic, have been demonstrated in the past to have catalytic properties for various types of hydrocarbon conversions. Certain zeolitic materials are ordered, porous crystalline aluminosilicates having a definite crystalline structure within which there are a large number of smaller cavities which may be interconnected by a number of still smaller channels. Since the dimensions of these pores are such as to accept for adsorption molecules of certain dimensions while rejecting those of larger dimensions, these materials have come to be known as \"molecular sieves\" and are utilized in a variety of ways to take advantage of these properties.\nSuch molecular sieves, both natural and synthetic, include a wide variety of positive ion-containing crystalline aluminosilicates. These aluminosilicates can be described as a rigid three-dimensional framework of SiO.sub.4 and AlO.sub.4 in which the tetrahedra are cross-linked by the sharing of oxygen atoms whereby the ratio of the total aluminum and silicon atoms to oxygen is 1:2. The electrovalence of the tetrahedra containing aluminum is balanced by the inclusion in the crystal of a cation, for example, an alkali metal or an alkaline earth metal cation. This can be expressed wherein the ratio of aluminum to the number of various cations, such as (Ca/2), (Sr/2), Na, K or Li is equal to unity. One type of cation may be exchanged either entirely or partially by another type of cation utilizing ion exchange techniques in a conventional manner. By means of such cation exchange, it has been possible to vary the properties of a given aluminosilicate by suitable selection of the cation. The spaces between the tetrahedra are occupied by molecules of water prior to dehydration.\nPrior art techniques have resulted in the formation of a great variety of synthetic aluminosilicates. These aluminosilicates have come to be designated by letter or other convenient symbols, as illustrated by zeolite A (U.S. Pat. No. 2,882-243), zeolite X (U.S. Pat. No. 2,882,244), zeolite Y (U.S. Pat. No. 3,130,007), zeolite ZK-5 (U.S. Pat. 3,247,195), zeolite ZK-4 (U.S. Pat. No. 3,314,752) and zeolite ZSM-5 (U.S. Pat. No. 3,702,886) merely to name a few.\nApplicants know of no prior art methods of zeolite preparation utilizing the present improvement. In fact, the present improved method of zeolite preparation is distinctly different from the current practice in the synthesis of zeolites having a silica/alumina mole ratio greater than 10."}
{"text": "1. Field of the Invention\nThis invention relates to computer memory systems and, more particularly, to techniques for avoiding data corruption during sparing of memory chips.\n2. Description of the Related Art\nEvolving standards of reliability, availability, and serviceability for computer memory systems have led to systems that include one or more redundant memory chips. In such systems, when failures from a particular memory chip exceed a threshold, a spare chip may be employed to take the place of the failed chip. This strategy, referred to as chip sparing, allows a multi-chip system to tolerate more than one failed memory chip by sparing out one known bad chip with a high number of failures and using error-correcting code (ECC) checking and correction to carry on without interruption when failures occur on other chips.\nA sparing operation typically includes sequentially reading the contents of a failed chip, applying error correction, and writing the corrected values to the spare chip. A simple approach to preventing data corruption is to stop other memory transactions during a sparing operation. However, stopping other memory transactions during a sparing operation may adversely affect system performance since chip sparing may take a significant amount of time to write-back the corrected contents of an entire chip. For example, to spare a DDR-533 512 Mb DRAM may take approximately 4 seconds. On the other hand, allowing a mixture of chip sparing transactions and other memory transactions could lead to data corruption. For example, when a non-sparing memory write to a given location occurs in between a read and a write of a chip sparing read/write pair to the same location (i.e. after the read but before the write), the data in the given location may be corrupted because the chip sparing write operation may overwrite the data written by the non-sparing write operation. Also, when a non-sparing memory write is issued to a given location that has been spared, data may be written to both the given location and to a corresponding spared location. In addition, a non-sparing memory read operation that takes place during a chip sparing operation to the same location may return different data depending on whether the read takes place before or after the chip sparing write-back. Accordingly, a system and method are needed that avoid data corruption during chip sparing operations while maintaining adequate system performance."}
{"text": "I. Field of the Invention\nThe present invention relates generally to orthodontic appliances and, more particularly, to an orthodontic appliance for imposing an anterior traction of force on the patient's jaw.\nII. Description of the Prior Art\nIn the practice of orthodontics, it is frequently necessary to move the patient's jaw anteriorly with respect to the patient's head in order to obtain a correct bite. This is especially true with the lower jaw of the patient.\nThere are a number of previously known appliances with are designed to move the patient's jaw forwardly by imposing anterior force on the patient's jaw. All of these previously known devices generally comprise a head gear or head assembly of one sort which is mounted to the patient's head so that a portion of the head gear is positioned in front of the patient's mouth. Thereafter, elastic members, typically rubber bands, extend between the appliance and the patient's jaw under tension in order to exert the anterior force on the patient's jaw.\nAs is well known in physics, for each force or reaction there must be an equal and opposite force or reaction. Consequently, in order to impose an anterior force on the patient's jaw, there must be an equal and opposite posterior force also imposed upon the patient.\nMany of these previously known devices have thus used pads which engage either the cheekbone area or the chin of the patient. Consequently, when rubber bands are tensioned between the previously known orthodontic device and the patient's jaw, a posterior force is imposed by the appliance against either the patient's cheekbone area, the patient's chin, or both. The rearward or posterior force imposed upon the cheekbone or chin is opposite to and offsets the anterior force imposed on the patient's jaw.\nHowever, both the chin as well as the cheekbones are interconnected with the patient's upper and lower jaws. Consequently, by imposing a posterior force on the cheekbones and/or the patient's chin, these previously known orthodontic appliances disadvantageously counteract the desired anterior force on the patient's jaw. Furthermore, in some cases, the posterior force imposed on the patient's chin and/or cheekbones result in some depression of the cheekbones and/or chin which detracts from the overall facial appearance of the patient."}
{"text": "Signaling through receptor tyrosine kinases (RTKs) regulates and fine-tunes many processes including cell growth, proliferation, differentiation, and apoptosis. The improper activation of RTKs is involved in the pathogenesis, growth, and metastasis of many cancers. The receptor tyrosine kinase ALK (Anaplastic Lymphoma Kinase) is a member of the insulin receptor superfamily that was initially identified from the t(2;5)(p23;q35) translocation in anaplastic large cell lymphoma (ALCL) (Fischer, P., et al. Blood, 72: 234-240. (1988)). The protein product of this translocation is ALK fused to nucleophosmin (NPM) (Morris et al., 1994). When fused to ALK, the dimerization domain of NPM results in constitutive dimerization and activation of ALK (reviewed in Chiarle, R., Nature reviews, 8:11-23 (2008)). Once activated, ALK recruits several adaptor proteins and stimulates multiple signaling pathways known to mediate tumor cell growth and survival including STAT3, PLC-γ, RAS-ERK1,2, and PI3K-AKT (Bai, R. Y., et al. Molecular and cellular biology 18: 6951-6961 (1998); Bai, R. Y., et al. Blood 96:4319-4327 (2000); Chiarle, R., et al. Nature medicine 11:623-629 (2005); Pulford, K., et al. Journal of cellular physiology 199:330-358 (2004)). The dysregulation of ALK is highly oncogenic, as it is sufficient to induce cell transformation in a several immortalized cell lines (Bischof, D., et al. Molecular and cellular biology 17:2312-2325 (1997); Fujimoto, J., et al. Proceedings of the National Academy of Sciences of the United States of America 93: 4181-4186 (1996)) and to form tumors in animal models (Chiarle, R., et al. Blood 101: 1919-1927 (2003); Kuefer, M. U., et al. Blood 90: 2901-2910 (1997)). Moreover, NPM-ALK drives tumor formation, proliferation and survival in ALCL (reviewed in (Duyster, J., et al. Oncogene 20: 5623-5637 (2001)).\nMore recently, ALK translocations have been detected in ˜5% of non-small cell lung cancers (NSCLC). Similar to ALK translocations in ALCL, the fusion proteins in NSCLC display constitutive ALK activity and drive tumor growth and survival (Soda et al., Nature 448: 561-566 (2007); Soda et al., Proceedings of the National Academy of Sciences of the United States of America 105: 19893-19897 (2008)). NSCLC tumors harboring ALK translocations are mutually exclusive from K-Ras or EGFR aberrations and predominantly occur in younger patients that are non-smokers (Rodig et al., Clin Cancer Res 15: 5216-5223 (2009); Shaw et al., J Clin Oncol 27: 4247-4253 (2009); Wong et al., Cancer 115: 1723-1733 (2009)). In addition to chromosomal rearrangements, activating point mutations and amplifications have been reported in a subset of sporadic and familial neuroblastomas, further expanding the spectrum of tumors dependent on ALK activity (Chen et al., Nature 455: 971-974 (2008); George et al., Nature 455: 975-978 (2008); Janoueix-Lerosey et al., Nature 455: 967-970 (2008); Mosse et al., Nature 455: 930-935 (2008)). Neuroblastomas with ALK genetic aberrations also are dependent on ALK for proliferation and survival, and cells expressing ALK containing activating mutations form tumors in animal models.\nInhibitors of RTKs have the potential to cause lethality in cancerous cells that are reliant on deregulated RTK activity while sparing normal tissues. Thus, small molecule inhibitors of ALK would be beneficial for therapeutic intervention in ALCL, NSCLC, neuroblastoma, and other cancers that are dependent on ALK for growth and survival."}
{"text": "As typified by prior art U.S. Pat. No.: 2,057,142 (Fry--Oct. 13, 1936); U.S. Pat. No. 2,738,580 (Rice--Mar. 20, 1956); U.S. Pat. No. 2,855,668 (Ottenad--Oct. 14, 1958); and U.S. Pat. No. 3,830,310 (Williams--Aug. 20, 1974): plunger-actuated surface-treating devices extending along an upright longitudinally extending central-axis and having a downwardly thrustable plunger member toward a lower-end surface-impacting tool are suggested in prior art teachings. However, such prior art plunger-actuated surface-treating devices are fraught with several disadvantages and deficiencies such as, but not limited to: are not structurally adaptable to efficiently dislodge variously randomly encountered topical-surfaces from their underlying substrates; are operationally unreliable for their intended purposes; are too heavy to be conveniently and reliably operated by variously physically endowed persons for the surface-treating intended purposes; and are not economically acquirable nor maintainable by the intended operational population."}
{"text": "1. Field of the Invention\nThe invention relates to the filtration of molten metal and, more particularly, to a relatively inexpensive molten metal filter that provides high flow rates while attaining excellent filtration capabilities.\n2. Description of the Prior Art\nIn the course of casting molten metal, the metal first is reduced to a molten state at an elevated temperature by means of a furnace. After being liquified in the furnace, the molten metal is directed into a mold for casting, or it is directed into a secondary receiver such as a holding furnace from which it later is directed into a mold for casting.\nOne of the major detriments that an occur during the casting process is for impurities in the form of solids to be included in the molten metal as it is poured into the casting molds. These impurities can originate from metal being melted itself, the furnace in which the metal is melted, transport elements through which the metal is conveyed, or external sources such as oxidation from the atmosphere. The inclusion of such impurities in the metal being cast can result in castings of poor quality that will require that the resultant castings be scrapped. In addition to problems associated with solid impurities, it sometimes happens that the molten metal contains dissolved gases, such as hydrogen, oxygen, or nitrogen. During solidification of the molten metal in the casting molds, the gases may precipitate out of solution and form gas bubbles or pockets that also may negatively affect the quality of the finished metal casting.\nIn view of the problems that can arise from the presence of solid or gaseous impurities in molten metal being cast, it is desirable to remove those impurities before solidification commences. A variety of filters have beer proposed for that purpose. In general, these filters employ a porous, refractory filter medium that will withstand the temperatures of the molten metal being filtered while, at the same time, trapping impurities that are being passed through the filter element. Such filters have been provided in the form of horizontally oriented plates, vertically oriented plates, and disks that are disposed within conduits through which the molten metal is flowed.\nA particularly effective molten metal filter is the MULTICAST filter system that is commercially available from the Metaullics Systems Division of The Carborundum Company, 31935 Aurora Road, Solon, Ohio 44139. In the MULTICAST filter system, a filter assembly is disposed within a holding furnace/filter box. The filter assembly includes a horizontally oriented sealing plate and several vertically oriented filter elements that are attached to the sealing plate. The filter elements ar&lt;hollow and surround openings in the sealing plate through which molten metal can flow. The filter elements are made of a porous ceramic material in the form of bonded silicon carbide or alumina particles that create numerous tortuous flowpaths that provide effective filtration capability. Reference is made to U.S. Patent Application Ser. No. 06/661,543, filed Oct. 16, 1984, by Edwin P. Stankiewicz, entitled \"Multiple-Use Molten Metal Filters,\", the disclosure of which is incorporated herein by reference.\nAnother type of molten metal filter is a so-called cartridge filter. In a cartridge filter, a plurality of horizontally oriented, refractory tubes are disposed within a holding furnace/filter box. The ends of the tubes are attached to, and supported by, graphite end plates. One of the end plates is solid, while the other end plate includes openings that are in fluid communication with the hollow interiors of the refractory tubes. Molten metal is filtered by disposing the cartridge filter within the holding furnace/filter box and by bringing molten metal into contact with the exterior surfaces of the tubes. The molten metal flows through the tubes (where it is filtered) and is discharged through the openings in the end plate for conveyance into the casting molds.\nDespite the filtration capabilities of prior art molten metal filters, particularly the MULTICAST molten metal filter and cartridge filters, certain problems remain. While the MULTICAST filter provides excellent filtration capability, its flow rate is less than desired for modern casting requirements. If molten aluminum is the metal being processed, the MULTICAST filter assembly can handle approximately 1,000-1,200 pounds per minute (with 6-grit particle size). Current industry trends demand that higher flow rates, i.e., 2,000-2,400 pounds per minute, be available. Desirably, an effective molten metal filter would be available that would have the filtration capability of the MULTICAST filter assembly, while providing higher flow rates.\nAlthough a cartridge filter provides acceptable flow rates, certain problems have not been addressed by conventional cartridge filters. In particular, a cartridge filter is large and expensive. The configuration of the filter elements is such that a great deal of wasted space is taken up by the filter. Also, the use of two large graphite end plates unduly increases the expense of the filter assembly. Desirably, a molten metal filter would be available that would (1) have a higher flow rate capability than the MULTICAST filter assembly, (2) have excellent filtration capability (as with the MULTICAST filter assembly), (3) occupy much less space than a cartridge filter, and (4) be much less expensive than a cartridge filter."}
{"text": "A 3D image system for displaying a 3D image, for example, an experience learning system projects a user's image in spaces including a subway station, a museum, and the like which are implemented virtually by using 3D and thereafter, enables a user's image object projected on a 3D content screen according to a user's actual action to visit the museum or learn English with a native speaker in line with a 3D content screen while performing a predetermined appointed action.\nIn this system, accuracy and stability of extracting a person image exert a great influence on total system performance. In the existing background subtraction based technologies, an object can be effectively extracted with a somewhat small calculation amount, but the system needs to be initialized whenever the system is driven and using a blue screen, and the like are disabled due to a feature of the system in which an experiencer's image needs to be extracted.\nIn the experience learning system, a learner that performs an experience and learners who perform learning while viewing the learner are present due to the feature thereof. In the existing object extracting system, it is possible to extract all users reflected on a viewing angle of a camera, but it is difficult to selectively extract a specific user."}
{"text": "Recent developments in the manufacture of thick blankets of intertangled mineral fibers include systems comprising a plurality, e.g. 4 to 6 or more, fiber forming and web forming modules, each producing a relatively thin, permeable, lightweight, web of intertangled mineral fibers. The webs from each of these modules are sequentially delivered to a transfer conveyor running under the modules to form a thick layer of superposed webs, which layer can later be processed into relatively thick mineral fiber blankets having a high degree of uniformity through the thickness of the blanket. These modular processes represented a significant improvement in the prior art technique of forming relatively thick mineral fiber blankets, particularly in the uniformity of the resultant product. Typical systems of the type described above are described in more detail in U.S. Pat. application Ser. No. 236,372 filed Mar. 20, 1972 and in U.S. Pat. No. 3,824,086, the disclosures of which are herein incorporated by reference.\nAlthough the above described systems represent a significant improvement in the manufacture of thick fiber glass blankets, there are features in these systems that need further improvement to permit the systems to be used more effectively. The initial capital expenditure required to build a system of the type described above, containing five or six modules, amounts to several millions of dollars. Once built, each module, in a pot and marble operation, requires an operator to insure that the fiberizing portion of the module is always operating properly. Although many products desired to be manufactured on these systems require that all of the modules operate simultaneously, many of the products are of a thickness that require the operation of less than all of the modules. When such products are being manufactured by the system described above, one or more modules remain idle, and the operators of these modules either remain idle or must be transferred temporarily to other duties, which is not always possible. Because of the large capital investment requirement of these systems, and operating costs that continue regardless of whether each module is operating or idle, the economics would be improved considerably if all of the modules could be operated continuously regardless of the requirements of the primary product being manufactured by the system.\nOne solution to this problem is discussed in Ser. No. 236,372 and amounts to dividing the transfer conveyer into two sections with the direction of travel of each sectional transfer conveyer being reversible. Then, if only a portion of the modules are required in the manufacture of the primary product, this portion of the modules can be selected from the modules at one end of the process, leaving the remaining modules on the opposite end of the system to manufacture a secondary product. The transfer conveyer under the modules not required for the primary product is reversed, these modules deliver their individual mats to the reversed transfer conveyer which delivers a secondary product of one, or two or more superposed layers, to the end of the system opposite the end to which the primary product is being delivered. While this modification can be used to solve the above described problem it adds to the building space requirements of the over-all system because it requires that additional processing equipment be located adjacent, and extending from, each end of the system. Also, additional personnel are required to operate the equipment at each end of the system simultaneously and this additional personnel becomes surplus when the primary product requires operation of all of the modules. Thus, this solution to the problem, while operable, is not entirely satisfactory from an economical standpoint.\nIt is an object of the present invention to modify the systems described above in such a manner that all or some of the modules can be operated independently of the overall system to manufacture separate and useful products when simultaneous operation of all of the modules are not required in the manufacture of the primary product.\nIt is a further object of the present invention to modify the systems described above to accomplish the above described objective using a minimum amount of additional equipment located in available space, i.e. without requiring additional building space requirements.\nIt is a further objective of the present invention to accomplish the described objectives and further to manufacture additional useful products without requiring additional operating personnel.\nIt is a further objective of the present invention to accomplish the above described objectives and also to provide the above described systems with increased production capacity and versatility."}
{"text": "Heretofore in dry wall construction systems, it has been the normal practice for a carpenter to secure or fasten the dry wall to the channel or standard framing members utilizing sheet metal screws and to thereafter have a mason secure the inner end of wall-ties to the previously secured dry wall and framing member assembly. It will therefore be apparent that there are two fastening operations required by two different trades to accomplish the preliminary securements prior to the construction of the outer wythe. This method is time consuming and thus more costly than one wherein a singular fastening operation can be utilized and performed."}
{"text": "1. Field of the Invention\nThis invention relates to a tamper protector enclosure assembly and is more particularly concerned with providing a tamper protector enclosure for an interface device within a distribution box and the associated telephone wiring or power line wiring feeding into and out of such interface device in order to discourage tampering and/or interruption of telephone service or electrical service through the interface device.\n2. Background Information\nTelephone network interface devices are well known to those skilled in the art, and typically comprise distribution box in which the telephone company's network wiring and the subscriber's wiring are joined at an interface device within the distribution box, as well as providing a demarcation point with regard to the ownership and maintenance responsibility for those telephone lines. Telephone network interface devices typically are provided with some means of sheltering the telephone line terminals joining the telephone company's and subscribers telephone lines. This shelter is usually a distribution box mounted externally on a house or building. The interface device within the box is usually easily accessible for determining whether problems in telephone service are the result of the telephone company's network, or the subscriber's wiring and equipment within the premises being served by the telephone company. The typical telephone network interface device, through which telephone service is provided to the subscriber's building is attached to the exterior of a structure, more typically on a residential rather than a commercial structure, but such devices are also used on commercial structures.\nIn the prior art, a number of telephone network interface devices exist which pertain and are primarily concerned with the manner and method in which the telephone company's lines are connected Co the subscriber's telephone service lines. The prior art boxes for the network interface devices generally make provision for enclosing the network interface device within a weather-resistant closure, capable of being locked. In a typical residential application, however, the distribution box and the network interface device are mounted upon the exterior of the residence, the cover for the box being only a plastic housing which fastens to the base which carries the network interface device, which, in turn, fastened to the structure.\nTypical of the prior art telephone network interface device and distribution box discussed, is an apparatus similar to that disclosed in U.S. Pat. No. 4,853,960 issued to Thomas J. Smith on Aug. 1, 1989 which provides an apparatus for connecting the telephone company's lines to the subscribers' service lines on a network interface baseplate, the baseplate having a covered housing which may be locked into a closed position. The device of Smith may be used for one pair of lines, or multiple pairs in a multi-unit residential or commercial setting.\nAnother type of network interface device and box disclosed in the prior art is found in U.S. Pat. No. 4,825,466, issued to Dowler et al. on Apr. 25, 1989 which provides a means for interfacing and enclosing a multiple number of telephone lines pairs, typically in a commercial or multi-unit residential facility. Another type of enclosure for a telephone network interface device is disclosed in U.S. Pat. No. 4,823,381 issued to Robert D. Olson on Apr. 15, 1989, which device comprises a locking security cover with an open frame for multiple telephone line pairs installed and connected to outside telephone line service upon a backboard. As with the Dowler patent, the Olson patent is intended primarily for use in either commercial or multi-unit residential facilities, and, as is the case with all the prior art, designed to protect the network interface only, and to limit access to the wired connections within the box of the interface structure.\nNone of these devices, however, relate to an apparatus which will completely enclose the interface device itself, as well as the associated telephone line wire from the telephone company, and to the subscriber, which feed into and exit the interface device. This is particularly important in those situations where the telephone company subscriber has installed in their residence, or in their business, a centrally wired burglar alarm system which is, itself, able to communicate with a central monitoring or control station through automatic dialing means over the subscribers' telephone lines and the telephone company's network lines.\nWhere a burglar is intent upon unlawfully entering a person's residence, or business, and is aware of the fact that the premises are guarded by a centrally wired burglar alarm system, it is the typical practice for the burglar to disable the central monitoring feature of the burglar alarm by cutting either or both of the telephone lines feeding into an interface device. Once this has been accomplished, the burglar knows that he will be able to gain ready access to the structure without fear of actuating the alarm system for contacting the central monitoring or control station, which would then notify the local law enforcement authorities to investigate the alarm signaled into the monitoring station.\nNone of the prior art recited above has taken into account or provided a means for enclosing the network interface device, the telephone company's service line, and the subscribers telephone line, feeding into and out of the interface device. Accordingly, even if a residential or commercial telephone subscriber has an interface device which has a locking cover, telephone service through said interface can be readily interrupted by severing the telephone lines leading into the interface. Also, none of the prior art has provided a means for signaling an alarm within the structure (distribution box) which contains the interface, nor to a central monitoring or control station. The preferred embodiment described, and as claimed, below, makes provision for fully enclosing the telephone network interface device, the subscriber's telephone lines and telephone company's telephone lines to thereby discourage the deliberate interruption of service, yet provide a means for ready access to the enclosed interface device and telephone line pairs to authorized telephone company and burglar alarm system repairmen.\nElectrical power lines are also mounted externally on a structure and provided with distribution boxes similar to those of the telephone service, the function being to house the interface network for the incoming power line and the outgoing power lines. The terminals and exposed wires at the interface are not, however, well protected by the distribution box from weather and these terminals and wires become corroded. This is particularly severe when the structure is close to salt water. The invention disclosed herein contemplates providing additional weather proofing protection for electrical power line distribution or junction boxes and the interfaces therein."}
{"text": "Diabetes refers to a disease process derived from multiple causative factors and characterized by elevated levels of plasma glucose or hyperglycemia in the fasting state or after administration of glucose during an oral glucose tolerance test. Persistent or uncontrolled hyperglycemia is associated with increased and premature morbidity and mortality. Often abnormal glucose homeostasis is associated both directly and indirectly with alterations of the lipid, lipoprotein and apolipoprotein metabolism and other metabolic and hemodynamic disease. Therefore patients with Type 2 diabetes mellitus are at especially increased risk of macrovascular and microvascular complications, including coronary heart disease, stroke, peripheral vascular disease, hypertension, nephropathy, neuropathy, and retinopathy. Therefore, therapeutical control of glucose homeostasis, lipid metabolism and hypertension are critically important in the clinical management and treatment of diabetes mellitus.\nThere are two generally recognized forms of diabetes. In Type 1 diabetes, or insulin-dependent diabetes mellitus (IDDM), patients produce little or no insulin, the hormone which regulates glucose utilization. In Type 2 diabetes, or noninsulin dependent diabetes mellitus (NIDDM), patients often have plasma insulin levels that are the same or even elevated compared to nondiabetic subjects; however, these patients have developed a resistance to the insulin stimulating effect on glucose and lipid metabolism in the main insulin-sensitive tissues, which are muscle, liver and adipose tissues, and the plasma insulin levels, while elevated, are insufficient to overcome the pronounced insulin resistance.\nInsulin resistance is not primarily due to a diminished number of insulin receptors but to a post-insulin receptor binding defect that is not yet understood. This resistance to insulin responsiveness results in insufficient insulin activation of glucose uptake, oxidation and storage in muscle and inadequate insulin repression of lipolysis in adipose tissue and of glucose production and secretion in the liver.\nThe available treatments for Type 2 diabetes, which have not changed substantially in many years, have recognized limitations. While physical exercise and reductions in dietary intake of calories will dramatically improve the diabetic condition, compliance with this treatment is very poor because of well-entrenched sedentary lifestyles and excess food consumption, especially of foods containing high amounts of saturated fat. Increasing the plasma level of insulin by administration of sulfonylureas (e.g. tolbutamide and glipizide) or meglitinide, which stimulate the pancreatic β cells to secrete more insulin, and/or by injection of insulin when sulfonylureas or meglitinide become ineffective, can result in insulin concentrations high enough to stimulate the very insulin-resistant tissues. However, dangerously low levels of plasma glucose can result from administration of insulin or insulin secretagogues (sulfonylureas or meglitinide), and an increased level of insulin resistance due to the even higher plasma insulin levels can occur. The biguanides increase insulin sensitivity resulting in some correction of hyperglycemia. However, the two biguanides, phenformin and metformin, can induce lactic acidosis and nausea/diarrhea. Metformin has fewer side effects than phenformin and is often prescribed for the treatment of Type 2 diabetes.\nThe glitazones (i.e. 5-benzylthiazolidine-2,4-diones) are a more recently described class of compounds with potential for ameliorating many symptoms of Type 2 diabetes. These agents substantially increase insulin sensitivity in muscle, liver and adipose tissue in several animal models of Type 2 diabetes resulting in partial or complete correction of the elevated plasma levels of glucose without occurrence of hypoglycemia. The glitazones that are currently marketed are agonists of the peroxisome proliferator activated receptor (PPAR), primarily the PPAR-gamma subtype. PPAR-gamma agonism is generally believed to be responsible for the improved insulin sensititization that is observed with the glitazones. Newer PPAR agonists that are being tested for treatment of Type II diabetes are agonists of the alpha, gamma or delta subtype, or a combination of these, and in many cases are chemically different from the glitazones (i.e., they are not thiazolidinediones). Serious side effects (e.g. liver toxicity) have occurred with some of the glitazones, such as troglitazone.\nAdditional methods of treating the disease are still under investigation. New biochemical approaches that have been recently introduced or are still under development include treatment with alpha-glucosidase inhibitors (e.g. acarbose) and protein tyrosine phosphatase-1B (PTP-1B) inhibitors.\nCompounds that are inhibitors of the dipeptidyl peptidase-IV (“DPP-4”) enzyme are also under investigation as drugs that may be useful in the treatment of diabetes, and particularly Type 2 diabetes. See WO 97/40832; WO 98/19998; U.S. Pat. No. 5,939,560; U.S. Pat. No. 6,303,661; U.S. Pat. No. 6,699,871; U.S. Pat. No. 6,166,063; Bioorg. Med. Chem. Lett., 6: 1163-1166 (1996); Bioorg. Med. Chem. Lett., 6: 2745-2748 (1996); Ann E. Weber, J. Med. Chem., 47: 4135-4141 (2004); D. Kim, et al., J. Med. Chem., 48: 141-151 (2005); and K. Augustyns, Exp. Opin. Ther. Patents, 15: 1387-1407 (2005). The usefulness of DPP-4 inhibitors in the treatment of Type 2 diabetes is based on the fact that DPP-4 in vivo readily inactivates glucagon like peptide-1 (GLP-1) and gastric inhibitory peptide (GIP). GLP-1 and GIP are incretins and are produced when food is consumed. The incretins stimulate production of insulin. Inhibition of DPP-4 leads to decreased inactivation of the incretins, and this in turn results in increased effectiveness of the incretins in stimulating production of insulin by the pancreas. DPP-4 inhibition therefore results in an increased level of serum insulin. Advantageously, since the incretins are produced by the body only when food is consumed, DPP-4 inhibition is not expected to increase the level of insulin at inappropriate times, such as between meals, which can lead to excessively low blood sugar (hypoglycemia). Inhibition of DPP-4 is therefore expected to increase insulin without increasing the risk of hypoglycemia, which is a dangerous side effect associated with the use of insulin secretagogues.\nDPP-4 inhibitors also have other therapeutic utilities, as discussed herein. DPP-4 inhibitors have not been studied extensively to date, especially for utilities other than diabetes. New compounds are needed so that improved DPP-4 inhibitors can be found for the treatment of diabetes and potentially other diseases and conditions. In particular, there is a need for DPP-4 inhibitors that are selective over other members of the family of serine peptidases that includes quiescent cell proline dipeptidase (QPP), DPP8, and DPP9 (see G. Lankas, et al., “Dipeptidyl Peptidase-IV Inhibition for the Treatment of Type 2 Diabetes,” Diabetes, 54: 2988-2994 (2005). The therapeutic potential of DPP-4 inhibitors for the treatment of Type 2 diabetes is discussed by D. J. Drucker in Exp. Opin. Invest. Drugs, 12: 87-100 (2003); by K. Augustyns, et al., in Exp. Opin. Ther. Patents, 13: 499-510 (2003); by J. J. Holst, Exp. Opin. Emerg. Drugs, 9: 155-166 (2004); by H.-U. Demuth in Biochim. Biophys. Acta, 1751: 33-44 (2005); by R. Mentlein, Exp. Opin. Invest. Drugs, 14: 57-64 (2005)"}
{"text": "1. Field of the Invention\nThis invention relates to coprecipitated nickel-cobalt-silica and nickel-cobalt-copper-silica hyrogenation catalysts and their use in hydrogenating organic compounds. In one aspect, this invention relates to the preparation of supported, coprecipitated nickel-cobalt-silica and nickel-cobalt-copper-silica catalysts. In another aspect, this invention relates to the use of the catalysts, in reduced form, to hydrogenate organic compounds.\n2. Description of the Prior Art\nThe catalytic reduction of organic compounds in the presence of nickel and cobalt as well as nickel-cobalt or nickel-cobalt-copper catalysts is known. Nickel catalysts, especially supported nickel catalysts, have many commercial uses. For example, U.S. Pat. No. 3,535,271 teaches the use of a nickel catalyst promoted by copper for dehydrogenation, cracking, reforming, polymerization, isomerization, alkylation, as well as other treating processes. Other examples of nickel catalysts and their use in reforming and other processes include U.S. Pat. Nos. 2,750,261; 3,205,182; 3,351,566; 3,417,029; 3,697,445; 3,859,370; 3,868,332; 4,088,603; and Belgium Pat. No. 841,812 (which generally corresponds to U.S. application Ser. No. 577,328). In all of these patents, the catalysts are prepared by coprecipitation or impregnation processes wherein the catalytic metal precursors are either precipitated from solution in the presence of a support material or solution containing said precursor or impregnated into the pores of a porous support material. In British Pat. No. 1,220,105, for example, aqueous solutions are employed in conjunction with a homogeneous precipitation procedure to give highly dispersed nickel catalyst.\nD. J. C. Yates, W. F. Taylor and J. H. Sinfelt (J. Am. Chem. Soc., 86, 2996 (1964)) described a chemisorption technique and its utility in correlating nickel particle size (and/or nickel surface area) with catalytic activity. In FIG. 3 of their publication, there is shown that a direct relation exists between reduced nickel surface area (m.sup.2 /g of catalyst) and initial reaction rate for ethane catalytically converted into methane (as mmoles C.sub.2 H.sub.6 converted per hour per gram of catalyst). It follows, then, that methods which increase the nickel surface area of a nickel catalyst (other factors such as nickel content remaining constant) is a desirable feature, leading to a catalyst of improved catalytic activity. Patentees of U.S. Pat. Nos. 3,697,445; 3,859,370 and 3,868,332 also appreciated that achieving a higher degree of dispersion of nickel in the catalyst results in a more active catalyst and indeed they obtain a fairly high degree of dispersion by their coprecipitation techniques wherein nickel cations were gradually precipitated from an aqueous solution in the presence of silicate anion and solid porous particles to obtain dispersion greater than 70 m.sup.2 /g of reduced nickel metal per gram of catalyst. Belgium Pat. No. 841,812 teaches that the addition of copper ions during the precipitation step provides a catalyst that can be reduced at temperatures of approximately 200.degree. C. U.S. Pat. No. 4,088,603 discloses an improved method of activating the coprecipitated nickel-copper-silica catalysts by gradually contacting the catalyst with a hydrocarbon feed and hydrogen.\nA number of patents have disclosed cobalt, nickel-cobalt and nickel-cobalt-copper catalysts, e.g., U.S. Pat. Nos. 3,166,491; 3,385,670; 3,432,443; 3,547,830; 3,650,713; 3,661,798; 3,945,944; 4,014,933 and 4,026,823; and British Pat. Nos. 1,000,828; 1,000,829; 1,095,996; 1,095,987 and 1,182,829. None of these patents, however, disclose coprecipitation of silicate ions in conjunction with nickel and cobalt or nickel, cobalt and copper ions in the presence of solid porous carrier particles.\nIn some of the above-mentioned patents, for example, U.S. Pat. Nos. 3,697,445; 3,859,370; 3,868,332 and Belgium Pat. No. 841,812 it is mentioned that cobalt or iron may be used in place of nickel in the coprecipitation process. However, these patents only show nickel in the examples as the non-noble Group VIII catalytic metal used."}
{"text": "1. Field of the Invention\nThe present invention relates to a segmented brake light, and more particularly to a brake light composed of multiple light modules that are selectively and progressively activated according to different pressure a driver stepped on the brake pedal.\n2. Description of Related Art\nA top-installed rear brake light near the rear windshield in a vehicle is usually started when the driver is stepping on a brake pedal to warn other people behind the vehicle to avoid any collision.\nHow the top-installed rear brake light works is often determined by the sensitivity of the brake light and a driver's behaviors. A high sensitivity light can be easily activated with a small pressure when the driver slightly steps on the brake pedal. However, the driver may just move the feet to the brake pedal as a preparatory action for braking, not mean to brake the car. Comparing to the high sensitive operations, a low sensitively brake light requires a deep stepping or high pressure on the brake pedal. For other drivers behind the car, they may have no sufficient time in emergency brakes. The complete turn-off and turn-on operations of a conventional light are unable to provide other drivers with a good pre-warning effect.\nHowever, to produce a top-installed rear brake light that can display multiple effects according to different statuses, some complex factors such as modifying original circuits in the vehicle may discourage consumers from equipping the vehicle with such a brake light.\nTherefore, the present invention provides a new segmented brake light to overcome the foregoing drawbacks."}
{"text": "Throughout this application, various publications are referenced by arabic numerals within parentheses. Full citations for these references may be found at the end of the specification immediately preceding the claims. The disclosures of these publications in their entireties are hereby incorporated by reference into this application in order to more fully describe the state of art as known to those skilled therein as of the date of the invention described and claimed herein.\nSuperoxide dismutase (SOD) and the phenomenon of oxygen free radicals (O2) was discovered in 1968 by McCord and Fridovich (1). Superoxide radicals and other highly reactive oxygen species are produced in every respiring cell as by-products of oxidative damage to a wide variety of macromolecules and cellular components (for review see 2,3). A group of metalloproteins known as superoxide dismutases catalyze the oxidation-reduction reaction 202+2H+xe2x86x92H2O2+O2 and thus provide a defense mechanism against oxygen toxicity.\nThere are several known forms of SOD containing different metals and different proteins. Metals present in SOD include iron, manganese, copper and zinc. All of the known forms of SOD catalyze the same reaction. These enzymes are found in several evolutionary groups. Superoxide dismutases containing iron are found primarily in prokaryotic cells. Superoxide dismutases containing copper and zinc has been found in virtually all eukaryotic organisms (4). Superoxide dismutases containing manganese have been found in organisms ranging from microorganisms to man.\nSince every biological macromolecule can serve as a target for the damaging action of the abundant superoxide radical, interest has evolved in the therapeutic potential of SOD. The scientific literature suggests that SOD may be useful in a wide range of clinical applications. These include prevention of oncogenesis and of tumor promotion, and reduction of the cytotoxic and cardiotoxic effects of anticancer drugs (10), protection of ischemic tissues (12) and protection of spermatozoa (13). In addition, there is interest in studying the effect of SOD on the aging process (14).\nThe exploration of the therapeutic potential of human SOD has been limited mainly due to its limited availability.\nSuperoxide dismutase is also of interest because of its anti-inflammatory properties (11). Bovine-derived superoxide dismutase (orgotein) has been recognized to possess anti-inflammatory properties and is currently marketed in parts of Europe as a human pharmaceutical.\nIt is also sold in the United States as a veterinary product, particularly for the treatment of inflamed tendons in horses. However, supplies of orgotein are limited. Prior techniques involving recovery from bovine or other animal cells have serious limitations and the orgotein so obtained may produce allergic reactions in humans because of its non-human origin.\nCopper zinc superoxide dismutase (CuZn SOD) is the most studied and best characterized of the various forms of superoxide dismutase.\nHuman CuZn SOD is a dimeric metallic-protein composed of identical non-covalently linked subunits, each having a molecular weight of 16,000 daltons and containing one atom of copper and one of zinc (5). Each subunit is composed of 153 amino acids whose sequence has been established (6,7).\nThe cDNA encoding human CuZn superoxide dismutase has been cloned (8). The complete sequence of the cloned DNA has also been determined (9). Moreover, expression vectors containing DNA encoding superoxide dismutase for the production and recovery of superoxide dismutase in bacteria have been described (24,25). The expression of a superoxide dismutase DNA and the production of SOD in yeast has also been disclosed (26).\nRecently, the CuZn SOD gene locus on human chromosome 21 has been characterized (27) and recent developments relating to CuZn superoxide dismutase have been summarized (28).\nMuch less is known about manganese superoxide dismutase (MnSOD). The MnSOD of E. coli K-12 has recently been cloned and mapped (22). Barra et al. disclose a 196 amino acid sequence for the MnSOD polypeptide isolated from human liver cells (19). Prior art disclosures differ, however, concerning the structure of the MnSOD molecule, particularly whether it has two or four identical polypeptide subunits (19,23). It is clear, however, that the MnSOD polypeptide and the CuZn SOD polypeptide are not homologous (19). The amino acid sequence homologies of MnSODs and FeSOD from various sources have also been compared (18).\nBaret et al. disclose in a rat model that the half life of human MnSOD is substantially longer than the half-life of human copper SOD; they also disclose that in the rat model, human MnSOD and rat copper SOD are not effective as anti-inflammatory agents whereas bovine copper SOD and human copper SOD are fully active (20).\nMcCord et al. disclose that naturally occurring human manganese superoxide dismutase protects human phagocytosing polymorphonuclear (PMN) leukocytes from superoxide free radicals better than bovine or porcine CuZn superoxide dismutase in xe2x80x9cin vitroxe2x80x9d tests (21).\nThe present invention concerns the preparation of a cDNA molecule encoding the human manganese superoxide dismutase polypeptide or an analog or mutant thereof. It is also directed to inserting this cDNA into efficient bacterial expression vectors, to producing human MnSOD polypeptide, analog, mutant and enzyme in bacteria, to recovering the bacterially produced human MnSOD polypeptide, analog, mutant or enzyme. This invention is also directed to the human MnSOD polypeptides, analogs, or mutants thereof so recovered and their uses.\nThis invention further provides a method for producing enzymatically active human MnSOD in bacteria, as well as a method for recovering and purifying such enzymatically active human MnSOD.\nThe present invention also relates to a DNA molecule encoding the human MnSOD gene. It is also directed to inserting the DNA into mammalian cells to produce MnSOD polypeptide, analog, mutant and enzyme.\nThe present invention also relates to using human manganese superoxide dismutase or analogs or mutants thereof to catalyze the reduction of superoxide radicals to hydrogen peroxide and molecular oxygen. In particular, the present invention concerns using bacterially produced MnSOD or analogs or mutants thereof to reduce reperfusion injury following ischemia and prolong the survival period of excised isolated organs. It also concerns the use of bacterially produced MnSOD or analogs thereof to treat inflammations.\nA DNA molecule which includes cDNA encoding the human manganese superoxide dismutase polypeptide or analog or mutant thereof has been isolated from a human T-cell library. The nucleotide sequence of a double-stranded DNA molecule which encodes human manganese superoxide dismutase polypeptide or analog or mutant thereof has been discovered. The sequence of one strand encoding the polypeptide or analog thereof is shown in FIG. 1 from nucleotide 115 downstream to nucleotide 708 inclusive. Other sequences encoding the analog or mutant may be substantially similar to the strand encoding the polypeptide. The nucleotide sequence of one strand of a double stranded DNA molecule which encodes a twenty-four (24) amino acid prepeptide is also shown in FIG. 1, from nucleotides number 43 through 114, inclusive.\nThe double-stranded cDNA molecule or any other double-stranded DNA molecule which contains a nucleotide strand having the sequence encoding the human manganese superoxide dismutase polypeptide or analog or mutant thereof may be incorporated into a cloning vehicle such as a plasmid or virus. Either DNA molecule may be introduced into a cell, either procaryotic, e.g., bacterial, or eukaryotic, e.g., yeast or mammalian, using known methods, including but not limited to methods involving cloning vehicles containing either molecule.\nPreferably the cDNA or DNA encoding the human manganese superoxide dismutase polypeptide or analog or mutant thereof is incorporated into a plasmid, e.g., pMSE-4 or pMSxcex94RB4, and then introduced into a suitable host cell where the DNA can be expressed and the human manganese superoxide dismutase (hMnSOD) polypeptide or analog or mutant thereof produced. Preferred host cells include Escherichia coli, in particular E. coli A4255 and E. coli A1645. The plasmid pMSE-4 in E. coli strain A4255 has been deposited with the American Type Culture Collection under ATCC Accession No. 53250. The plasmid pMS RB4 may be obtained as shown in FIG. 4 and described in the Description of the Figures.\nCells into which such DNA molecules have been introduced may be cultured or grown in accordance with methods known to those skilled in the art under suitable conditions permitting transcription of the DNA into mRNA and expression of the mRNA as protein. The resulting manganese superoxide dismutase protein may then be recovered.\nVeterinary and pharmaceutical compositions containing human MnSOD or analogs or mutants thereof and suitable carriers may also be prepared. This human manganese superoxide dismutase or analogs or mutants may be used to catalyze the following reaction:\n202+2H+xe2x86x92H2O2+O2\nand thereby reduce cell injury caused by superoxide radicals.\nMore particularly, these enzymes or analogs or mutants thereof may be used to reduce injury caused by reperfusion following ischemia, increase the survival time of excised isolated organs, or treat inflammations.\nThis invention is directed to a method of producing enzymatically active human manganese superoxide dismutase or an analog or mutant thereof in a bacterial cell. The bacterial cell contains and is capable of expressing a DNA sequence encoding the manganese superoxide dismutase or analog or mutant thereof. The method comprises maintaining the bacterial cell under suitable conditions and in a suitable production medium. The production medium is supplemented with an amount of Mn++ so that the concentration of Mn++ available to the cell in the medium is greater than about 2 ppm.\nIn a preferred embodiment of the invention the bacterial cell is an Escherichia coli cell containing a plasmid which contains a DNA sequence encoding for the human manganese superoxide dismutase polypeptide e.g. pMSE-4 or pMSxcex94RB4 in E. coli strain A4255. The concentration of Mn++ in the production medium ranges from about 50 to about 1500 ppm, with concentrations of 150 and 750 ppm being preferred.\nThis invention also concerns a method of recovering manganese superoxide dismutase or analog thereof from bacterial cells which contain the same. The cells are first treated to recover a protein fraction containing proteins present in the cells including human manganese superoxide dismutase or analog or mutant thereof and then the protein fraction is treated to recover human manganese superoxide dismutase or analog or mutant thereof. In a preferred embodiment of the invention, the cells are first treated to separate soluble proteins from insoluble proteins and cell wall debris and the soluble proteins are recovered. The soluble proteins are then treated to separate, e.g. precipitate, a fraction of the soluble proteins containing the hMnSOD or analog or mutant thereof and the fraction containing the hMnSOD or analog or mutant is recovered. The recovered fraction of soluble proteins is then treated to separately recover the human manganese superoxide dismutase or analog thereof.\nA more preferred embodiment of the invention concerns a method of recovering human manganese superoxide dismutase or analog or mutant thereof from bacterial cells which contain human manganese superoxide dismutase or analog or mutant thereof. The method involves first isolating the bacterial cells from the production medium and suspending them in suitable solution having a pH of about 7.0 to 8.0. The cells are then disrupted and centrifuged and the resulting supernatant is heated for about 30 to 120 minutes at a temperature between 55 and 65xc2x0 C., preferably for 45-75 minutes at 58-62xc2x0 C. and more preferably for 1 hour at 60xc2x0 C. and then cooled to below 10xc2x0 C., preferably to 4xc2x0 C. Any precipitate which forms is to be removed e.g. by centrifugation, and the cooled supernatant is dialyzed against an appropriate buffer e.g. 2 mM potassium phosphate buffer having a pH of about 7.8. Preferably, the dialysis is by ultrafiltration using a filtration membrane smaller than 30K. Simultaneously with or after dialysis the cooled supernatant optionally may be concentrated to an appropriate convenient volume e.g. 0.03 of its original volume. The retentate is then eluted on an anion exchange chromatography column with an appropriate buffered solution e.g. a solution of at least 20 mM potassium phosphate buffer having a pH of about 7.8. The fractions of eluent containing superoxide dismutase are collected, pooled and dialyzed against about 40 mM potassium acetate, pH 5.5. The dialyzed pooled fractions are then eluted through a cation exchange chromatography column having a linear gradient of about 40 to about 200 mM potassium acetate and a pH of 5.5. The peak fractions containing the superoxide dismutase are collected and pooled. Optionally the pooled peak fractions may then be dialyzed against an appropriate solution e.g. water or a buffer solution of about 10 mM potassium phosphate buffer having a pH of about 7.8.\nThe invention also concerns purified enzymatically active human manganese superoxide dismutase or analogs thereof e.g. met-hMnSOD, or mutants produced by the methods of this invention.\nThe present invention also relates to a DNA molecule encoding the human MnSOD gene. The nucleotide sequence of the exon coding regions of one strand of the MnSOD gene is shown in 6. The DNA encoding the MnSOD gene may be incorporated into a cloning vehicle such as a plasmid or a virus. The DNA or the cloning vehicle may be introduced into eucaryotic cells using known methods. Preferably, the DNA encoding the human manganese superoxide dismutase gene is encoded in the plasmids pMSG11-1, pMSG4 and pMSG-1b. The eucaryotic cells which are transformed are preferably mammalian cell lines such as the human HeLa cell line or the mouse L cell line. Another aspect of this invention is the production of the human manganese superoxide dismutase polypeptide, analog, mutant or enzyme by growing the cells of this invention in suitable medium under suitable conditions."}
{"text": "Scintillation detectors have been employed in various industrial applications, such as the oil and gas industry for well logging. Typically, such detectors have scintillator crystals made of an activated sodium iodide material that is effective for detecting gamma rays. Generally, the scintillator crystals are enclosed in tubes or casings, which include a window permitting radiation induced scintillation light to pass out of the scintillator casing for measurement by a light-sensing device such as a photomultiplier tube. The photomultiplier tube converts the light photons emitted from the scintillator into electrical pulses that are shaped and digitized by associated electronics that may be registered as counts and transmitted to analyzing equipment. In terms of well logging applications, the ability to detect gamma rays makes it possible to analyze rock strata as gamma rays are emitted from naturally occurring radioisotopes, typically of shales that surround hydrocarbon reservoirs.\nDesired properties of scintillation detectors include high resolution, high brightness, and robust design to withstand mechanical and thermal loads, while providing a high quality hermetic seal to isolate the sensitive internals of the detector from harsh operating environments. In this latter aspect, a common practice in the utilization of scintillation detectors is to take measurements while drilling (MWD). For MWD applications, the detector should be ruggedized, that is, it should have a high shock resistance and be capable of withstanding high temperatures, while maintaining performance specifications for a reasonable lifetime.\nRare-earth halides are known to have improved scintillation properties such as resolution and brightness. However, such materials have been difficult to integrate into detectors rugged enough to withstand harsh environments and high temperatures.\nAs such, there continues to be a need for improved scintillation detectors, particularly scintillation detectors incorporating rare-earth halides or other state of the art scintillator materials that can withstand the rigors of industrial applications."}
{"text": "1. Technical Field\nThe present disclosure relates to methods and valves for controlling the flow of fluid through a bore and more particularly, the disclosure relates in some embodiments to methods and ball valves for use in the oil and chemical process industry. More particularly, this disclosure relates to double piston, trunnion-mounted ball valves.\n2. Background Art\nBall valves are commonly used in both the oil and chemical process industries. A type of ball valve used to control flow of a fluid is an apertured ball valve such as is disclosed in PCT Patent Application No. WO 93/03255 published on Feb. 18, 1993, incorporated by reference herein. In an apertured ball valve the valve operation or function may be broken down into two separate stages. Firstly, the ball moves between an open and a closed position by rotating through 90 degrees such that the ball aperture from an orientation coaxial with the flow direction, i.e. when the valve is open, to a position whereby the ball aperture is normal or perpendicular to the flow direction and the valve is closed. Secondly, the valve seals in the closed position to prevent flow through the bore across the ball valve. Therefore, the on-off control of flow through the valve is achieved by rotating the ball through 90° within the valve housing. Another ball valve is disclosed in U.S. Pat. No. 6,708,946, the teachings of which are incorporated herein by reference.\nThere are two basic types of ball valve mechanisms which currently exist. First, there is the trunnion mounted ball system in which the ball element is positionally constrained inside the valve, usually by radial bearings. The ball is rotated by the application of torque through a valve stem to the trunnion. Sealing occurs as a result of the valve seat on the upstream (or high pressure) side of the valve “floating” onto the ball element and causing engagement between a surface of the valve seat and the surface of the ball. The advantage of this system is that it provides highly reliable rotation between the valve open and the closed positions. The principal disadvantage of this system is that seal reliability is reduced because the sealing force only develops in proportion to the annular area of the valve seat. Further, in high pressure applications, the force exerted on the ball on the upstream side of the valve can result in deformation of the ball and leakage between the ball and another valve seat located on the downstream (or lower pressure) side of the valve. Thus, when trunnion mounted ball systems are used in high pressure wells and especially those in which the well fluid has a high proportion of particulate matter, being generally known as “aggressive” wells, the pressure is such that fluids and/or particulate matter may leak past seals between the ball and the valve seats. This often results in the valve not achieving integrity of sealing. In such cases, this type of ball valve is unable to operate properly in such conditions.\nThe second type of ball valve mechanism which effects the abovementioned function is known as the “floating ball system”. In this system the ball is not positionally constrained relative to the valve body. Rotation is caused by the application of force to a point which is offset from the ball centre which, in conjunction with the mating curvatures of the ball and seat, cause the ball to rotate. Sealing occurs as a result of the ball “floating” onto the valve seat. The advantage of this mechanism is that the reliability of the seal is increased, because the sealing force develops in proportion to the circular area of the ball to seat contact. The disadvantage of this type of mechanism is that the rotational reliability is reduced as the friction factor between the ball and seat are considerably larger than that of trunnion mounted devices. With high pressure and aggressive types of wells and particulate flows of the type described above, the reliability of this valve in those applications creates a problem in that the torque necessary to rotate the ball becomes excessively high, and thus, the valve can seize between the open and the closed position giving rise to serious problems in both operational and safety terms.\nIt would be desirable to provide a method and/or improved ball valve design which may obviate or mitigate at least one or more of the aspects associated with the aforementioned disadvantages."}
{"text": "This application is an application under 35 U.S.C. Section 371 of International Application Number PCT/FR98/02897 filed on Dec. 29, 1998.\nThe present invention relates to a method for preparing polyorganosiloxanes (silicones) with an ester function (acetoxyalkyl or acetoxyalkenyl) of high purity, primarily free of functions or by-products with an xcex1-olefinic unsaturated bond.\nThe invention provides a method for producing polyorganosiloxanes with ester functions with formula (I):\nR1R2R3SiO(R4R5SiO)p(R6QSiO)qSiR3R2R1xe2x80x83xe2x80x83(I),\nwhere:\nsymbols R1, R2 are identical or different and represent:\na C1 to C6 alkyl radical or a phenyl radical, preferably methyl;\nor a xe2x80x94OR7 radical, where R7 represents a linear or branched C1 to C6 alkyl radical, preferably methyl;\nsymbols R4, R5 and R6 are identical or different and represent a C1 to C6 alkyl radical or a phenyl radical, preferably methyl;\nsymbols R3, which may be identical or different, represent:\na linear or branched C1 to C6 alkyl radical or a phenyl radical, preferably methyl;\nthe symbol Q;\nor a xe2x80x94OR7 group, where R7 represents a linear or branched C1 to C6 alkyl radical, preferably methyl;\nsymbols Q, which may be identical or different, represent an ester function or an alcohol function, respectively with formulae (1) and (2) below:\nxe2x80x94Rxe2x80x2OCORxe2x80x3xe2x80x83xe2x80x83(1)\nxe2x80x94Rxe2x80x2OHxe2x80x83xe2x80x83(2)\nwhere:\nRxe2x80x2 represents a linear or branched C3-C20 alkylene group, preferably C3-C12, more particularly trimethylene or 2-methyltrimethylene;\nRxe2x80x3 represents a methyl group;\nat least one of symbols Q representing an ester function with formula (1);\np represents an average number of 0 or more, preferably more than 5;\nq represents an average number of 0 or more, preferably more than 1; at least one of symbols R3 representing the ester symbol Q of formula (1) when q equals 0;\nby hydrosilylation between a hydrogeno-polyorganosiloxane with formula II):\nR1R2R3SiO(R4R5SiO)p(R6HSiO)qSiR3R2R1xe2x80x83xe2x80x83(II)\nwhere:\nsymbols R1, R2, R4, R5, R6, p and q have the definitions given above;\nsymbols R3xe2x80x2, which may be identical or different, represent:\na linear or branched C1 to C6 alkyl radical or a phenyl radical, preferably methyl;\nH;\nor a xe2x80x94OR7 group, where R7 represents a linear or branched C, to C6 alkyl radical, preferably methyl;\nand an xcex1-olefinic alcohol with formula ROH, where R represents an xcex1-olefinic radical from which the divalent radical Rxe2x80x2 derives by opening of the xcex1 double bond and adding hydrogen, preferably an allyl or methallyl radical;\nthen esterification of the hydroxy-alkylated polyorganosiloxane with formula (III) obtained:\nR1R2R3xe2x80x3SiO(R4R5SiO)p(R6Qxe2x80x2SiO)qSiR3xe2x80x3R2R1xe2x80x83xe2x80x83(III)\nwhere:\nsymbols R1, R2, R4, R5, R6, p and q have the definitions given above;\nsymbols R3xe2x80x3, which may be identical or different, represent:\na linear or branched C1 to C6 alkyl radical or a phenyl radical, preferably methyl;\nthe symbol Qxe2x80x2;\nor a xe2x80x94OR7 group, where R7 has the definition given above;\nsymbol Qxe2x80x2 represents the radical Rxe2x80x2OH, where Rxe2x80x2 has the definition given above; with an acyl halide with formula XCORxe2x80x3, where X represents a halogen atom, in particular chlorine, and Rxe2x80x3 has the definition given above, in the absence of an esterification catalyst;\nand eliminating the xcex1 olefinic by-products and the hydrohalic acid formed.\nThe hydrosilylation step can be carried out in any known manner at a temperature of the order of 20xc2x0 C. to 200xc2x0 C., preferably of the order of 60xc2x0 C. to 120xc2x0 C., in the presence of a hydrosilylation catalyst. Examples which can be cited are those based on platinum such as the platinum complexes described in U.S. Pat. No. 3,159,601, U.S. Pat. No. 3,159,662, U.S. Pat. No. 3,715,334, U.S. Pat. No. 3 ,14,730. A conventional catalyst is the KARSTEDT catalyst used, for example, in an amount of 1 to 300 parts, preferably 5 to 50 parts, by weight of platinum per million parts by weight of reactants.\nThe relative quantities of hydrogeno-polyorganosiloxane with formula (II) and xcex1-olefinic alcohol ROH generally correspond to an xcex1-olefinic alcohol/hydrogeno-polyorganosiloxane mole ratio more than 1 and less than or equal to 5, preferably more than 1 and less than or equal to 2.\nIn defining the moles of (xcex1-olefinic alcohol, the elementary entity is considered to be the xcex1-olefinic unsaturated bond; in defining the moles of hydrogeno-polyorganosiloxane, the elementary entity is considered to be the SiH function.\nThe hydrosilylation operation can be carried out at atmospheric pressure, preferably in bulk; however, it can be carried out in an organic solvent such as toluene, xylene, methylcycyclohexane, heptane, octane, hexamethylsiloxane.\nThe operation is carried out until the xe2x80x94SiH groups in the reaction mass have disappeared. The operation can be monitored by infrared analysis or volumetric determination.\nThe excess xcex1-olefinic alcohol can be eliminated by vacuum distillation (for example 13 Pa to 14000 Pa).\nThe hydrogeno-polyorganosiloxane used in the hydrosilylation step is preferably selected from those with formulae:\nMe3SiO(SiMe2O)pxe2x80x2(SiMeHO)qxe2x80x2SiMe3\nMe2HSiO(SiMe2O)pxe2x80x2(SiMeHO)qxe2x80x2SiHMe2\nMe2HSiO(SiMe2O)pxe2x80x3SiHMe2\nwhere:\nMe represents the methyl radical;\npxe2x80x2 represents an average number of the order of 0 to 500, preferably of the order of 5 to 100;\npxe2x80x3 represents an average number of the order of 0 to 100, preferably of the order of 0 to 20;\nqxe2x80x2 represents an average number of the order of 1 to 50, preferably of the order of 1 to 30.\nThe xcex1-olefinic alcohol with formula ROH used in the hydrosilylation step is preferably allyl alcohol or methallyl alcohol.\nThe esterification operation can be carried out at a temperature of the order of 0xc2x0 C. to 100xc2x0 C., preferably of the order of 15xc2x0 C. to 50xc2x0 C., more particularly of the order of 15xc2x0 C. to 30xc2x0 C.\nThe relative quantities of hydroxyalkylated polyorganosiloxane and acyl halide with formula XCORxe2x80x3 generally correspond to an acyl halide/hydroxyalkylated polyorganosiloxane mole ratio of the order of 0.1 to 2, preferably of the order of 0.1 to 1.\nIn defining the moles of hydroxyalkylated polyorganosiloxane, the elementary entity is considered to be the hydroxyalkyl function.\nThe esterification operation can be carried out at atmospheric pressure, preferably in bulk; however, it can be carried out in an organic solvent such as toluene, xylene, methylcyclohexane, heptane, octane, hexamethylsiloxane.\nThe acyl halide with formula XCORxe2x80x3 used in the esterification step is preferably acetyl chloride.\nThe acid halide formed by esterification can be eliminated by entraining it in an inert gas (argon, nitrogen . . . ) then by treatment with a base, such as alkali metal bicarbonates, in particular sodium bicarbonate, in stoichiometric excess over the acyl halide (for example 2 to 3 times the stoichiometry).\nAfter eliminating the by-products by condensing under reduced pressure, the polyorganosiloxane with an ester function can then be recovered by filtering, for example.\nxcex1-olefinic siloxane motifs and xcex1-olefinic impurities are those with formulae xe2x95x90Sixe2x80x94Oxe2x80x94R, ROH and Rxe2x80x3COOR, where R represents a linear or branched xcex1-olefinic radical containing 3 to 20 carbon atoms, as defined above (allyl or methallyl, for example), and Rxe2x80x3 is the methyl group."}
{"text": "There is an ever increasing demand for low cost, high volume solutions for lens assemblies with adjustable focal length. Modern mobile phones, for example, are now equipped with miniature digital camera modules and the quality and cost demands for lenses and lens assemblies are increasing. More and more miniature cameras used in for example mobile phones and laptop computers comes with auto focus functionality. The design of lens systems for such applications requires fulfilment of a number of requirements, such as ease of operation when fitting the lens on top of a camera module, by providing as few operational steps as possible, etc. These challenges are even greater when the lens arrangement comprises tuneable parameters, such as encountered in auto focus lenses, wherein the focal length must be adjusted to fit the distance from the lens to the object to be photographed, for example. Such lenses are usually complex designs comprising movable parts which can make it difficult to assemble the lens in a simple manner. A further challenge with such designs is the ever increasing requirement to provide as thin lens assemblies as possible. Thin and lightweight mobile phones and cameras are a must in the marketplace.\nThere exist a number of feasible solutions for making compact auto focus lens elements. For example, liquid filled lens elements based on the principle of electro wetting is disclosed in WO2006136613. Another example is disclosed in the article “Design, fabrication and testing of a micro machines integrated tuneable micro lens”, by Weisong et. Al., Journal of Micromechanics and Microengineering, 9 May, 2006, wherein a lens assembly containing a liquid in a cavity may be compressed or elongated to achieve an adjustment of focal length of the lens.\nBoth these two examples of prior art lenses comprise one or more liquids. Filling and confining a liquid to a cavity may impose problems during both production and use, as known to a person skilled in the art.\nAnother principle for a tuneable lens assembly is described in U.S. Pat. No. 4,802,746, wherein a plurality of elastic members is superposed in the optical axial direction. By applying a force on a side of the lens element, a deformation creates a curvature that results in variable focal properties, depending on the magnitude of the applied force.\nJP 02178602 disclose a lens with an adjustable focal length comprising a transparent fluid 4 bounded by two transparent surfaces 1d, 2 and an elastic sealing material 3. The curvature of one of the transparent surfaces (1d) is adjusted by applying voltages onto a transparent piezo electric film 1c disposed onto the transparent surface.\nJP 01140118 disclose a lens with an adjustable focal length comprising two transparent surfaces 1,3 with a fluid 2 in between bounded by a wall 10. The fluid 2 can be a gel or polymer. One of the transparent surfaces is a transparent piezo electric polymer, wherein the curvature of this surface is controlled by applied voltages onto the two electrodes 4, 5.\nUS 2002/0048096 disclose how to control deformations of a mirror 1. Piezo electric actuators 3 are arranged in a pattern on the backside of the mirror, and the deformation is controlled based on signals received from sensors attached to the mirror.\nThe Norwegian patent application No. 20065237 disclose a lens assembly comprising a floating glass cover on top of a gel or elastomer disposed on a transparent support. Piezoelectric actuators on top of the floating glass cover located in the periphery of the floating glass cover provides bending of the floating glass cover when activated that results in a change of the surface curvature of the gel or elastomer. However, the range of focal length adjustments is limited due to the limited range of bending of the floating glass cover. If the actuating force is to high, the edge of the floating glass cover may cut the gel or elastomer located directly under the edge of the floating glass cover thereby damaging the lens body, and hence the range of focal length adjustments are limited by this possibility of damaging the lens body.\nTherefore there is a need to provide a lens assembly that provides the principal benefits of a lens assembly as disclosed in the Norwegian patent application No. 20065337, and that at the same time provides an increased range of focal length adjustments without changing properties of the lens assembly when adjustments are performed, and that protect the lens body against any form of damage when adjustment is performed in the whole range of adjustments, and that protect the lens body from damage or interference during manufacturing, when making interconnections, etc. of other optical and electronic components used together with the lens assembly according to the present invention. This property of the lens assembly according to the present invention also protects the lens assembly when using the lens assembly in an application."}
{"text": "1. Field of the Invention\nThe invention is concerned with a portable apparatus for conveniently, safely and accurately weighing infants and measuring their length.\n2. Prior Art\nU.S. Pat. No. 945,530--K. B. HAYDEN discloses an infant's weighing device comprising a blank of pliable material stitched on corded lines to impart a corrugated appearance and a cup-like form. Also included is a drawstring to further shape the device to a cup-like structure and to be able to adjust the size to some extent. Extensions of the cord are provided with rings to permit attachment to a weighing device. This device is uncomfortable for an infant, handling the infant was difficult and dangerous.\nU.S. Pat. No. 1,252,824--A. MELNIKER discloses a foldable cradle including foldable side frame members, collapsible end frames and means for hanging the cradle from suspending hooks.\nU.S. Pat. No. 1,853,198--C. E. BREADEN discloses a weighing scale platform for packages which also incorporates measuring means along three axii for measuring the size of a package.\nU.S. Pat. No. 440,959--P. A. O'MARRA discloses a combination desk ruler and paper or package weighing scale. The apparatus shown in this patent is awkward to use and uncomfortable for the infant. No means were provided for readily measuring the length and circumference of the infant nor is it constructed to fold down flat."}
{"text": "Immunoglobulin compositions prepared from human plasma and suitable for intravenous administration are known in the art and for several decades have played an important role in the treatment of a wide range of diseases. Immunoglobulins are used, for example, for the treatment of infections in humans and can be assigned to various classes with various biochemical and physiological properties. Immunoglobulin G participates in defending against viral antigens, whereas IgM is predominantly active in antibacterial and antitoxin immune responses.\nThe immunoglobulin solutions comprise IgG, IgA and IgM in various percentages, with different preparations having different treatment applications, e.g. preparations with a higher percentage of IgM are used in the prophylaxis or treatment of bacterial infections.\nThe immunoglobulin solutions are usually prepared from fractions of blood plasma or serum, e.g. Cohn fractions. These fractions are then subjected to a number of purification steps to remove contaminants such as viruses, denatured proteins, proteases and lipids.\nHuman plasma for fractionation is collected from thousands of donors and may contain pathogen viruses despite thorough testing of the source plasma. Therefore process steps to inactivate or remove viruses are essential in order to achieve safe products for use in medicine. Several techniques for virus inactivation/removal are known in the art, e.g. chemical treatments, irradiation with UVC light or nanometer filtration, which are performed in order to ensure overall virus safety.\nThe virus removal or inactivation capacity of the process steps is validated using laboratory scale models of the production process and for each step a removal or inactivation factor is determined. An increase of the inactivation/removal factor adds additional viral safety to the pharmaceutical product. Today guidelines from regulatory authorities require at least two effective steps for enveloped and non-enveloped viruses in the manufacture of plasma-derived pharmaceuticals. Although several methods, such as solvent/detergent treatment, octanoic acid treatment, nanometer filtration and heat treatment, are effective to inactivate or remove enveloped viruses there are only a few methods known to inactivate or remove non-enveloped viruses, for example Parvo viruses. These non-enveloped viruses are mostly very small, usually passing through nanometer filters with pore sizes above 20 nm. This pore size is too small for IgM molecules having a diameter up to 30 nm. Non enveloped viruses are effectively inactivated by chemicals like β-propiolactone which, however, also leads to a modified immunoglobulin with impaired functions. Another effective treatment is UVC-irradiation (EP1842561, CAF-DCF). However, known solvent/detergent treatments, octanoic acid treatment and mild heat treatment have no substantial effect on non-enveloped viruses.\nAs mentioned above, in addition to viruses which are potentially present it is also necessary to remove other contaminants like lipids, proteases, protein aggregates, and denatured immunoglobulins. The removal of all these contaminants is essential (1) to ensure the product complies with bio-safety guidelines regarding viral contamination, (2) in order for the product to be tolerated by the patient after intravenous administration, (3) to allow the product to be stable during long-term storage (any residual proteolytic activity might lead to degradation of the product over long-term storage, e.g. 2 years), and (4) to generate the desired compound mixture/pharmaceutical composition.\nAt the same time, however, it is essential that the purification steps to remove the contaminants do not interfere with the immunoglobulin molecules, so that as far as possible these retain their normal biological activity and are retained at high yield in solution. This balance is difficult to achieve since many known purification steps can also have a negative impact on the activity of the immunoglobulins, and in particular on IgM; for example extended irradiation times with UVC can reduce the yield of native and active IgM obtained in the final immunoglobulin solution. Not only does this lead to a reduction in efficacy of the final immunoglobulin solution but it can also cause the solution to be less well tolerated in vivo.\nAggregates and denatured immunoglobulins, the amount of which can be increased by certain purification steps, especially are a potential risk for the patients because they have a high capacity to activate complement unspecifically, leading to severe side effects in patients receiving these denatured immunoglobulins. Unspecific complement activation refers to the initiation of the complement cascade in the absence of specific antibody-antigen complexes. Unspecific complement activation is strictly to be avoided since it may cause undesirable side effects such as hypotension, flushing, headache, fever, chills, nausea, vomiting, muscle pain, dyspnoea and tachycardia. Specific complement activation, on the other hand, is desirable and it occurs only after the immunoglobulins have bound to their specific antigens.\nUnspecific complement activation is measured as the so called anticomplementary activity (ACA) by a standardized test described in the European Pharmacopoeia.\nThe role of the complement system in the immune defense of pathogens is well known. The complement system consists of about 20 proteins, which are activated sequentially. The classical complement pathway typically requires a specific antigen antibody complex for activation, whereas the alternative pathway can be activated by antigens without the presence of antibodies. The classical and the alternative pathway of complement activation all generate a protease C3-convertase. The C3-convertase cleaves and activates component C3, creating C3a and C3b, and causing a cascade of further cleavage and activation events over C5 convertase to C5a and C5b. C5b initiates the membrane attack pathway, which results in the membrane attack complex, consisting of C5b, C6, C7, C8, and polymeric C9. This is the cytolytic endproduct of the complement cascade which forms a transmembrane channel, which causes osmotic lysis of the target cells like bacteria.\nComplement activation additionally results in the formation of anaphylatoxins, including the biologically active protein C5a. This anaphylatoxin is a potent chemotactic agent for immune and inflammatory cells and induces cell activation and causing the release of histamine from mast cells. In situations of excessive or uncontrolled and/or unspecific complement activation, the overproduction of C5a can cause deleterious effects to patients.\nC5a is an effective leukocyte chemoattractant, causing the accumulation of white blood cells, especially neutrophil granulocytes, at sites of complement activation. C5a activates white blood cells and is a powerful inflammatory mediator. Whereas these functions are beneficial during specific antibody-antigen complex reactions all unspecific generation of C5a has to be avoided due to the potential side effects.\nUnspecific complement activation is a particular issue for IgM immunoglobulin preparations (i.e. those comprising at least 5% IgM) as in contrast to IgG preparations IgM antibodies easily aggregate in solution. IgM preparations are difficult to stabilize especially if they are enriched compared to plasma concentrations and stored in liquid solution. It is also known that IgM is a vigorous activator of complement; a single molecule bound to an antigen can activate complement. This is in contrast to IgG, where two or more molecules of IgG must be bound to an antigen in close association with each other to activate complement.\nStill further, the main indications treated by IgM containing immunoglobulin preparations are bacterial infections and sepsis. As these patients are already suffering from hypotension an additional unwanted generation of unspecific complement activation and C5a would lead to a clinical worsening of the patient's condition. Accordingly, IgM preparations have been described as being difficult to prepare for intravenous application.\nThere are several methods described in the art for the production of IgM containing immunoglobulin preparations from human plasma.\nThe initial purification of human IgM solutions has been carried out by classical Cohn plasma fractionation methods or its well known modifications (e.g. Cohn/Oncley, Kistler/Nitschmann). Using cold ethanol precipitation processes the IgM fraction is recovered in fraction III or fraction I/III (also called B or B+I). Starting from fraction III or I/III methods have been described for purification of protein solutions enriched in IgM. EP0013901 describes a purification method starting from fraction III including steps using octanoic acid, β-Propiolactone treatment and an adsorption step using an anionic exchange resin. This method is used to produce Pentaglobin®—to date the only commercially available intravenous IgM product. β-propiolactone is a well known chemical used in sterilization steps in order to inactivate viruses which are potentially present. As β-propiolactone is a very reactive substance which causes the chemical modification of proteins there is also substantial loss of the anti-viral and anti-bacterial activities of the immunoglobulins. On the other hand this chemical modification results in an reduced anticomplementary activity compared to an chemically unmodified immunoglobulin. EP0352500 describes the preparation of an IgM concentrate for intravenous application with a reduced anti-complementary activity by using anionic exchange chromatography, β-Propiolactone, UVC light irradiation and an incubation step at increased temperature (40° C. to 60° C.). The preparation produced by this method was stable in liquid solution for a limited time due to the chemical modification. The IgM concentration was above 50% from the total immunoglobulin content.\nThe preparation of protein solutions enriched in IgM without chemical modification by β-propiolactone has been described in EP0413187 (Biotest) and EP0413188 (Biotest). These methods involve subjecting a suitable protein solution to octanoic acid treatment and anionic exchange chromatography, starting from Cohn fraction III or II/III. In patent EP0413187 (Biotest) the octanoic acid treatment is carried out by stirring for 15 min, in order to remove lipids being present in Cohn fraction III.\nThe preparation according to EP0413187 had a low anticomplementary activity, between 0.6 and 0.8 CH50/mg protein, but had to be stabilized and virus inactivated by β-propiolactone. Low anticomplementary activity is considered to be ≦1 CH50/mg protein according to EP monograph for immunoglobulins.\nEP0413188B1 (Biotest) describes the preparation of an IgM-enriched preparation for intravenous administration by using an anion exchange chromatography in order to reduce the anti-complementary activity. Additionally a heat treatment at pH 4-4.5 at 40 to 60° C., preferably between 50 and 54° C., was described to reduce the anticomplementary activity. This preparation had to be lyophilized to ensure stability of the preparation for several months. Long term stability as a liquid solution could not be shown.\nM. Wickerhauser et al. “Large Scale Preparation of Macroglobulin”, Vox Sang 23, 119-125 (1972) showed that IgM preparations isolated by PEG precipitation had high anticomplementary activity (ACA) by a standard complement fixation test and this ACA activity was reduced 10 fold by incubating the IgM preparation at pH 4.0 at 37° C. for 8 hours followed by readjustment to neutral pH. It was not demonstrated if this 10 fold reduction is sufficient to ensure intravenous tolerability. The authors did not assess the specific complement activating potential of their IgM concentrate, nor did they assess safety in any animal or human model.\nAnother method describes the use of mild-heat treatment of IgM preparations at 40 to 62° C., preferably 45 to 55° C., at pH 4.0 to 5.0 (EP 0450412, Miles) to reduce the unspecific complement activation. In this patent application octanoic acid is added to a Cohn fraction III suspension in order to remove prekallikrein activator and lipoproteins by centrifugation. Nevertheless this mild heat treatment led to partial loss of antigenic determinants of IgM. This may increase the risk of generating neo-antigens leading to a increased immunogenicity in humans or the loss of activity.\nThe preparation of an IgM containing protein solution for intravenous application by using a protease treatment (e.g. with pepsin) after an octanoic acid precipitation step has been described in EP0835880 (U.S. Pat. No. 6,136,312, ZLB). Protease treatment leads to partial fragmentation of the immunoglobulin molecule impairing the full functional activity of the Fab and Fc parts. Therefore protease-treated immunoglobulins cannot be regarded as unmodified. Also this preparation method leads to about 5% fragments with a molecular weight of <100 kD.\nThe described methods to carry out the octanoic acid treatment (EP0413187 and EP0835880) have the drawback that the octanoic acid treatment is not effective with respect to removal and inactivation of non-enveloped viruses, and does not remove substantially all proteolytic activity.\nIn EP 0345543 (Bayer, Miles) a highly concentrated IgM preparation with at least 33% IgM for therapeutic use is disclosed, the preparation being substantially free of isoagglutinin titres. In this patent application an octanoic acid precipitation is carried out by adding the octanoic acid and the isoagglutinins are removed by Synsorb affinity chromatography. The final preparation had to be freeze dried.\nAltogether the production of an IgM containing preparation with low anticomplementary activity is possible if the immunoglobulins are chemically or enzymatically modified and/or further purified by chromatography and/or subjected to a mild heat treatment. However, these methods have their drawbacks in the lack of virus removal/virus inactivation (and therefore virus safety), reduction in the amount of immunoglobulin molecules in native form and/or residual anticomplementary activity. As such, there is still a need to provide improved IgM containing immunoglobulin preparations suitable for intravenous administration in humans."}
{"text": "1. Field of the Invention\nThe present invention is directed to technology for monitoring applications in a computing environment.\n2. Description of the Related Art\nThe growing presence of the Internet as well as other computer networks such as intranets and extranets has brought many new applications in e-commerce, education and other areas. Organizations increasingly rely on such applications to carry out their business or other objectives, and devote considerable resources to ensuring that the applications perform as expected. To this end, various application management techniques have been developed.\nFor example, for component-based applications, such as those based on JAVA™ 2 Platform Enterprise Edition (J2EE™) (a type of platform for server programming), one approach is to collect data about the individual software components that are invoked in an application. A software component generally refers to a software technology for encapsulating software functionality. A software component has characteristics including: multiple-use (reusable), non-context-specific, composable with other components, encapsulated, and a unit of independent deployment and versioning. Alternatively, a component can be considered to be an object written to a specification, such as JAVA BEANS™ (a type of software component), or the Component Object Model, a MICROSOFT platform for software componentry. A component can also include a web server or a database, for instance. Software components are typically invoked by processes or threads executing in an application, middleware or other software. For example, for a web-based e-commerce application, a process at an application server is typically initiated by a request to the server, such as a user request to purchase an item. The process may invoke a number of components to carry out the steps necessary to allow the user to purchase the item. For example, a shopping cart component may be used to allow the user to identify the item desired, the payment method and the shipping information. A reservation component may be used to reserve the item while a credit check component verifies the user's credit card information. Once the credit card information is verified, an inventory component is updated based on the item purchased, and a shipping component is invoked to arrange for the item to be shipped, such as by contacting a warehouse. An order completion component may be invoked to confirm the purchase to the user such as by providing an order confirmation number and a tracking number on a web page. Moreover, a given component may be invoked more than once during an application.\nIn particular, an execution flow can be traced to identify each component that is invoked as well as obtain performance data such as the execution time of each component. An execution flow refers generally to the sequence of steps taken when a computer program executes. Tracing refers to obtaining a detailed record, or trace, of the steps a computer program executes. One type of trace is a stack trace. Traces can be use as an aid in debugging. However, information cannot be obtained and analyzed from every execution flow without maintaining an excessive amount of overhead data and thereby impacting the very application which is being monitored. One way to address this problem is by sampling so that information is obtained regarding every nth execution flow. This approach is problematic because it omits a significant amount of data and, if a particular execution flow instance is not selected for sampling, all information about it is lost. Thus, if a particular component is executing unusually slowly, for instance, but only on an irregular basis, this information may not be captured.\nAnother approach, aggregation, involves combining information from all execution flows into a small enough data set that can be reported. For example, assume there are one thousand requests to an application server. For each execution flow, performance data such as the response time can be determined. Information such as the slowest, fastest, median and mean response times can then be determined for the aggregated execution flows. However, aggregating more detailed information about the execution flows is more problematic since the details of the execution flows can differ in various ways. Moreover, a mechanism has not been available for aggregating information between related execution flows, such as at different computer systems.\nIt would be desirable to provide a technique for monitoring execution flows which addresses the above and other issues."}
{"text": "1. Technical Field\nThe present invention relates in general to environmental sensing and control and in particular to a local processor using sensors for continuously monitoring environmental conditions within a data processing system. Still more particularly, the present invention relates to a local processor that utilizes registers arranged in a fault/mask/cache fashion to pass information to system firmware for environmental control and sensing within a data processing system or information handling system.\n2. Description of the Related Art\nMany data processing or computer systems support a standard input/output (I/O) systems conforming to the peripheral component interconnect (PCI) Local Bus architecture, an architecture supporting many complex features including I/O expansion through PCI-to-PCI bridges, peer-to-peer (device-to-device) data transfers, multi-function devices, and both integrated and plug-in devices. These input/output sub-systems may typically be set up in I/O drawer configurations, especially in large server systems having multiple I/O sub-systems. One of the complexities involved in these types of configurations is keeping the I/O drawer at the manufactures recommended operating temperatures and/or keeping enough supply power to run all the devices and operations. Even though personal computers or servers, during normal system operation, run little risk of corrupting data, the risk of data corruption becomes significant when environmental variables change (i.e. temperature) or system components (i.e. power supplies) become defective. This makes environmental sensing a very important feature.\nTherefore, it would be ideal if a computer system would monitor environmental data in the background and alert its system only when a change was sensed to protect itself from data corruption. However, a problem arises when trying to bridge the localized environmental sense information of a system component to the system level where an appropriate action can be taken due to hardware and software complexities. Consequently, it would be desirable to provide a method and system for monitoring and controlling at the I/O sub-system level environmental and system component information through an arrangement of simple hardware registers common to both the system firmware and the I/O drawer processor code. The present invention solves these problems in a novel and unique fashion not previously known in the art."}
{"text": "Remote memory, prepaid accounts for use in purchasing goods and services are generally well known. Presently known schemes typically include a printed transaction card, for example a wallet-sized plastic or cardboard card, which bears on one side a unique authorization or account number and instructions for access to funds, services, and the like. Such prepaid cards have been used extensively throughout the world. One such example is the use of these remote memory cards as prepaid long distance telephone calling cards.\nIn contrast to stored value cards (e.g., \"smart cards\") wherein a remaining account balance is stored within a memory resident in the card, remote memory systems generally store information pertaining to a prepaid account at a central host computer. The host computer typically stores information relating to the available balance remaining in the account, as well as information pertaining to past activity associated with the account. In particular, the host computer may store transaction data relating to various goods or services purchased using the card. In the context of a prepaid telephone calling card, the host computer may store telephone call data, including the date, time, duration, and various other parameters relating to calls which were placed using the prepaid telephone card. The host computer may be accessed via a telephone or data line by the consumer through the use of an authorization code, personal identification number (PIN), or the like.\nThe use of prepaid remote memory telephone cards is particularly prevalent in the United States. A typical prepaid telephone calling card includes a toll-free telephone number used by a consumer to access a host computer system, a unique authorization code associated with the card and a corresponding remote account, and dialing instructions. When a customer desires to use the card to place a long distance call, he dials the toll-free number, thereby accessing the host system which manages the remote accounts. The consumer then enters a predetermined authorization number for allowing access to the database. Next, the consumer enters the desired long distance telephone number, and the system connects the consumer with the desired calling destination. Long distance telephone charges attributed to the telephone call are deducted from the remaining balance in the account and the call is terminated when the account is fully consumed. A call history, which includes information pertaining to the calls charged to the various accounts, may also be maintained by the host computer for each account.\nIn the United States, preprinted instructions on how to use the transaction cards are generally printed directly on the cards in a single language, such as English. The use of such cards by individuals who are not conversant in the English language can be inconvenient or even impossible, depending on their fluency in the English language and their familiarity with how to use the transaction cards. While instructions may be printed on a single card in several languages, the number of languages that may be printed on a card is directly related to the size of the card and the level of detail to which the instructions will be printed; in the event wallet-sized cards are used, limited space usually precludes more than one language being printed on the card. Another option in attempting to resolve this language problem is to print a series of cards in a variety of different languages. However, this option has the disadvantage, among others, of increasing the issuer's and merchant's inventory by a factor equal to the number of languages printed. A third option is to assign each language a separate telephone number. Once connected to the host computer by phone, the consumer can be instructed on how to use the transaction card in that particular language. This however, has the disadvantage of requiring the issuer to maintain a phone number and connection for each language.\nConsequently, a system and method is thus needed which will overcome the shortcomings of the prior art and will allow the issuer to easily produce transaction cards that may be used by consumers in a plurality of different languages."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication system, a radio communication terminal and a radio communication device, for setting up a communication link while carrying out an authentication processing, which are directed to the case of carrying out transmission/reception of information by a short range radio communication such as Bluetooth between the radio communication terminal and the radio communication device, for example.\n2. Description of the Related Art\nRecently, the short range radio communication scheme having a transmission range of about 10 m has been attracting much attention. In contrast to the conventional radio LAN which has an effective transmission range of over 100 m, the short range radio communication scheme requires a smaller power consumption because of the short transmission range, so that it is suitable for an information processing device with a limited battery capacity such as a portable telephone.\nThe Bluetooth is one such short range radio communication scheme, and the portable information processing device (referred hereafter simply as a portable terminal) such as a portable telephone, a PDA, etc., which is equipped with a communication function according to this Bluetooth has been gaining increasingly popularity in recent years. For further details of the Bluetooth, the detailed specification can be obtained from the Web site “http://www.bluetooth.org/”.\nWhen the Bluetooth becomes widespread, it becomes possible to provide a service suitable for each location, at any location. For example, at the convenience store, supermarket or retail store, it is possible to realize services for the benefit of purchasers such as an electronic coupon service, an electronic point service, an electronic payment, an electronic receipt, an electronic acknowledgement issuance, etc. Furthermore, applications to a gate controlling using an electronic ticket, payment and discount at the automatic vending machine, an information delivery, a road guidance, etc. are expected.\nIn the following, an exemplary case where a connection according to the Bluetooth is set up between a portable terminal equipped with the communication function according to the Bluetooth (hereafter simply referred to as a Bluetooth portable terminal) and an information communication device equipped with the Bluetooth (hereafter referred to as a Bluetooth device), and the Bluetooth device provides a service to a user who owns the portable terminal will be described.\nThe Bluetooth uses a frequency band that is open to public for commercial use, so that it can be easily obstructed by other devices using the same frequency band. For this reason, the Bluetooth employs the frequency hopping technique for periodically changing the utilized frequency so as to be less vulnerable to the obstruction from the other devices. In the frequency hopping, there is a need for an initial operation to establish synchronization between two devices for carrying out communications.\nAt a time of searching the other Bluetooth devices existing in its communication area, the portable terminal uses the Bluetooth device search function to sequentially scan through the predetermined frequency band and discover the Bluetooth devices that are standing by in an accessible state.\nThe Bluetooth device in the accessible state that received a search signal transmits a clock information along with its own terminal identification information to a search signal sender. The Bluetooth device in the accessible state is also periodically changing the utilized frequency, so that it takes several seconds before the frequencies coincide. Even under the ideal condition of having no radio communication device using the same frequency in a vicinity, it may take five seconds or more to discover one terminal depending on the circumstance. In the case where the device using the same frequency exists in a vicinity, it can possibly take even longer time.\nThe portable terminal collects a terminal identification information in 48 bits from each one of the discovered Bluetooth devices, and carries out the synchronization of the frequency hopping pattern with a desired (selected) Bluetooth device at a connection phase by using the clock information and 28 bits out of the terminal identification information of that Bluetooth device. After the hopping pattern is synchronized, a physical link is established, and thereafter layers are connected so that the exchange and setting of necessary parameters are carried out, and then it becomes possible to carry out the data exchange at the application level. The mutual authentication is carried out if necessary, and when the connection at the application level is established, the Bluetooth device provides a prescribed service.\nIn general, in comparison with the communication utilizing infrared rays, the radio communication has an advantage in that it is not affected by the positional relationship or obstacles between two devices, but it also has a property that it is difficult to set up a connection only with a terminal existing at a specific location.\nThis fact gives rise to the following problems. In the following, an exemplary case where a user who has the electronic coupon in the portable terminal uses that electronic coupon by carrying out communications with a cash register equipped with the Bluetooth communication function at a time of payment at a cash register of the convenience store will be considered.\nIn the case where only one customer (only one portable terminal) and only one cash register (only one Bluetooth device) are present, the relationship between the portable terminal of the customer who makes the payment and the cash register is uniquely determined so that no problem regarding the connection arises. In other wards, when a person who operates the portable terminal discovers the Bluetooth device by using the Bluetooth device search function described above, the discovered device can be regarded as the cash register to which the payment is to be made.\nHowever, when a plurality of cash registers exist, the cash register of the next line will also be discovered, so that there is a need for the user who has the portable terminal to select the one cash register to which the payment is to be made from the discovered cash registers, by checking a name of the cash register or the like.\nAlso, in the case where the other customer waiting for his turn to make the payment also has a portable terminal, there is a possibility for the portable terminal of the other customer to make the connection with the cash register before the portable terminal of the customer currently in his turn to make the payment.\nFor this reason, there is a need for an operator of the cash register to check the customer if the connection is made with a right portable terminal or not, by using some information associated with the portable terminal.\nSuch selection and checking operations do not constitute a good interface for the customer who makes the payment and they also hinder the register operation of the operator.\nMoreover, if the person operating the portable terminal specifies the wrong cash register, there can be troubles such as the amount to be paid is not discounted despite of the use of the coupon, or the points granted to the one customer are actually added to the other customer. The cancellation operation must be carried out at both the portable terminal and the cash register, which requires considerable time and effort.\nIn addition, in the system using radio, anyone can be accessible to the radio communications, so that there can possibly be a problem related to business obstruction and privacy due to an attack from a malicious user such as setting up a link earlier than an innocent user to prevent the innocent user to make the connection, or pretending the cash register to secretly acquire the personal information or steal coupons, etc.\nThus, conventionally, it has been difficult to establish one-to-one correspondence between the portable terminal and the Bluetooth device correctly, easily and quickly, by using the short range radio communications."}
{"text": "Epigenetic regulation of gene expression is an important biological determinant of protein production and cellular differentiation and plays a significant pathogenic role in a number of human diseases.\nEpigenetic regulation involves heritable modification of genetic material without changing its nucleotide sequence. Typically, epigenetic regulation is mediated by selective and reversible modification (e.g., methylation) of DNA and proteins (e.g., histones) that control the conformational transition between transcriptionally active and inactive states of chromatin. These covalent modifications can be controlled by enzymes such as methyltransferases (e.g., PRMT5), many of which are associated with specific genetic alterations that can cause human disease.\nDisease-associated chromatin-modifying enzymes (e.g., PRMT5) play a role in diseases such as proliferative disorders, metabolic disorders, and blood disorders. Thus, there is a need for the development of small molecules that are capable of inhibiting the activity of PRMT5."}
{"text": "Urinary incontinence affects over 13 million men and women of all ages in the United States. Stress urinary incontinence (SUI) affects primarily women and is generally caused by two conditions, intrinsic sphincter deficiency (ISD) and hypermobility. These conditions may occur independently or in combination. In ISD, the urinary sphincter valve, located within the urethra, fails to close properly (coapt), causing urine to leak out of the urethra during stressful activity. Hypermobility is a condition in which the pelvic floor is distended, weakened, or damaged, causing the bladder neck and proximal urethra to rotate and descend in response to increases in intra-abdominal pressure (e.g., due to sneezing, coughing, straining, etc.). The result is that there is an insufficient response time to promote urethral closure and, consequently, urine leakage and/or flow results.\nA popular treatment of SUI is the use of a sling, which is permanently placed under a patient's bladder neck or mid-urethra to provide a urethral platform. Placement of the sling limits the endopelvic fascia drop, while providing compression to the urethral sphincter to improve coaptation. Generally, the sling is surgically placed under urethra, close to the high-pressure zone with no elevation of the urethra. The ends of a sling are secured to an anatomical location by, for example, a bone anchor. Recent advances in surgical techniques have demonstrated the effectiveness of approaches without employing bone anchors toward mid-urethra sling stabilization. However, the bone anchor-less techniques typically require the use of large diameter delivery devices. For example, some procedures require ⅛ to ¼ inch diameter delivery devices to create a tract in which the sling would lie. Such a large tract may cause the patients pain for weeks. Further, the typical sling made from woven mesh may also narrow when tensioned during placement, which results in less stable placement in tissue.\nAccordingly, there is a need for an improved surgically implantable sling or sling-like support having a low profile, and which does not narrow during placement."}
{"text": "Conventionally, the carrier device is known that carries a work such as a substrate in a vacuum space (for example, see patent document 1). The carrier device includes a linkage arm mechanism that horizontally moves the work in the vacuum space.\nIn addition, the radiation heat is conveyed to a peripheral member of the work because the work that is carried by the carrier device is in a high temperature state. Therefore, the carrier device is known that includes a rail mechanism and in which a cooling path is provided in a guide rail portion of the rail mechanism that is susceptible to the radiation heat from the work in the high temperature state (for example, see patent document 2)."}
{"text": "U.S. Pat. Nos. 4,024,175 and 4,087,544, which are herein incorporated by reference, disclose novel cyclic amino acids of Formula A ##STR1## wherein R.sub.1 is a hydrogen atom or a lower alkyl radical and n is 4, 5, or 6 and the pharmacologically compatible salts thereof.\nThe compounds disclosed in the above United States patents are useful for the therapy of certain cerebral diseases, for example, they can be used for the treatment of certain forms of epilepsy, faintness attacks, hypokinesia, and cranial traumas. Additionally, they bring about an improvement of cerebral functions and thus are useful in treating geriatric patients. Particularly valuable is 1-(aminomethyl)-cyclohexaneacetic acid (gabapentin).\nGamma-aminobutyric acid (GABA) is an inhibitory amino acid found in the mammalian central nervous system (CNS). It has been reported that dysfunction with GABA neurotransmission in the CNS may contribute or even cause psychiatric and neurological diseases, such as epilepsy, schizophrenia, Parkinson's disease, Huntington's Chorea, and dyskinesia (Saletu, B, et al, International Journal of Clinical Pharmacology, Therapy and Toxicology, 24, pages 362 to 373 (1986)). Gabapentin was designed as a GABA analog that would cross the blood-brain barrier. Gabapentin was found to have anticonvulsant and antispastic activity with extremely low toxicity in man.\nThe aforementioned compounds of Formula A including gabapentin have been prepared, from a compound of formula ##STR2## wherein R.sub.2 is an alkyl radical containing up to eight carbon atoms and n is as defined above, by well known standard reactions such as, for example, the Hofmann, Curtius, or Lossen rearrangements, into the amino derivatives of Formula A. Although these reactions provide the target compounds they require a large number of synthetic steps and in some cases involve potentially explosive intermediates.\nU.S. Pat. No. 4,152,326 discloses cyclic sulphonyloxyimides of formula ##STR3## wherein R.sub.2 is a saturated, straight-chained, branched or cyclic lower aliphatic radical or an unsubstituted or substituted aryl radical and n is 4, 5, or 6 which can be converted into a compound of Formula A. Again, similar to the previous processes, this process requires a large number of synthetic steps to obtain a compound of Formula A. Finally, all of the previous processes require as the penultimate step conversion of an intermediate salt of the target compound to an amino acid of Formula A.\nThe object of the present invention is an improved process for preparing the compounds described above by using a novel synthesis.\nFurther, we have unexpectedly found that the particularly valuable gabapentin can be prepared from a novel intermediate without the need for proceeding through an intermediate salt in the penultimate step and in fewer steps and higher yields than the previous methods. Moreover, the present method proceeds from inexpensive starting materials and is amenable to large-scale synthesis."}
{"text": "Optical connectors are used to connect optical signals between two optical components, such as between two optical fibers, between an optical fiber and a piece of optical or optoelectronic equipment, or between two pieces of optical or optoelectronic equipment. Due to its many advantages, optical fibers and other optical transports, such as optical waveguides, are increasingly being used to transport signals within individual pieces of electronic equipment, such as computers, cellular telephones, televisions, etc. Accordingly, optical connections must be made across printed circuit boards that are movable relative to each other. For instance, a removable daughter card, such as a PCMCIA card might mate to a motherboard through the back panel of a personal computer. There are several different relative orientations in which two PCB boards (or cards) may mate, including coplanar, orthogonal, and conventional backplane, each having a different relative orientation between the two mating PCBs.\nFurthermore, there are many different styles of optical connectors. Merely a few of the well-known standard optical connector styles are LC, MT, MPO, SC, and ST. Each connector style has its own form factor and can be mated only to another connector of the same style.\nAs data rates and bandwidths in electronic devices continue to increase, it sometimes is the case that a manufacturer of an electronic device wishes to upgrade the type of optical connector style used in a card without having to replace the entire card. For example, an LC connector can have up to four optical transports in it. If a daughter card, such as a PCMCIA card for a personal computer was originally designed to mate with a backplane using LC connectors, the manufacturer eventually may be interested in upgrading the optical couplings from LC style connectors to MT style connectors, which can support up to 72 fibers per connector, in order to increase the bandwidth of the data transfer between the motherboard and the daughter card.\nTypically, this requires a complete redesign of both the mother board and the daughter card."}
{"text": "Magnetic recording stores information in a hard magnetic layer of a magnetic medium by setting the direction of magnetization of regions of the layer. A recording apparatus of the type using a magnetic disk, magnetic tape or similar magnetic recording medium, reproduces stored digital data in the form of an analog waveform representative of transitions of magnetic polarity on the recording medium. In the case of data recorded in an NRZ (non-return to zero) or NRZI (non-return to zero inverted) recording pattern, for example, positive and negative peaks follow one another in succession with transitions between the peaks passing through the signal baseline. Depending on the code used, the positive and negative peaks may represent explicit data bit values, in self-clocking codes, or may represent explicit clock bits and encoded data bits in pulse position or interval modulation codes. One form of a digital data reproducing circuit for reading any such recorded code differentiates the reproduced waveform and determines the points where the differentiated waveform crosses a zero level representative of zero AC volts, to detect the peaks.\nIn commonly assigned U.S. Pat. No 4,977,419, to Wash et al., a photographic filmstrip having a virtually transparent, magnetic film layer covering the non-emulsion side of the filmstrip (referred to as an MOF layer) is disclosed in conjunction with various camera systems. Potential longitudinal recording tracks are illustrated in the MOF layer in both the image frame area and on either side of the image frame area for recording information such as film type, film speed, film exposure information and information relevant to the processing and subsequent use (e.g., printing) of the exposed image frames. The systems disclosed therein provide for recording of information during film manufacture, reading and/or recording of information on certain tracks during camera use, and reading and/or recording of printer related information during photofinishing.\nNovel methods for modulating the binary data or information into a format suitable for recording and reproducing in such tracks in the camera shown in the '419 patent are disclosed in commonly assigned U.S. Pat. No. 4,964,139, also to Wash et al. Two self synchronizing PPM coding methods are disclosed in the '139 patent, as the \"Wash encoded signal\" and the \"Chi encoded signal\" methods. Both PPM encoding methods have similarities as pointed out in the '139 patent, and the Wash encoded signal method is further referred to herein for purposes of explaining how the present invention may be practiced in at least one of its preferred embodiments in locating drop-outs in magnetically recorded data encoded in such a manner.\nAn example of the Wash encoding method, as used to record data in bit streams in the MOF layer tracks described in the '419 patent, is shown in FIG. 1 which essentially reproduces FIG. 3 of the '139 patent. In FIG. 1A), two successive bits 0 and 1 are position encoded in the two corresponding \"interval or information-cells\" as positive going data signal level transitions bounded by negative going clock signal transitions between the directions of magnetization. The Wash self-clocking bit code is effected by the position of the positive going data signal level transition in the information-cell. If it is positioned within the first half of the duration of the information-cell, then a data bit 0 is recorded. Similarly, a data bit 1 is recorded by positioning the data transition in the second half of the information-cell. As an aside, it should be noted that the data bit may in fact be a parity bit, as the two terms are used below.\nNote that the encoding of the information or data transitions for both the first and second information-cells leaves invariant the clock transitions. With this self-clocking code, film transport velocity can vary during recording and playback without affecting the ability to synchronize and read the recorded data. As a result, the information-cells may lengthen or shorten depending on the filmstrip velocity, but the relative position of the data transitions can still be decoded. Thus, the camera disclosed in the '419 patent, for example, may record data in the MOF layer tracks while winding the filmstrip between exposures without imposing any velocity controls or recording an independent clock track. In fact, large amounts of jitter in the camera transport at frequencies at or near the data frequencies precludes the use of a phase-locked loop in generating a clock during magnetic signal read out.\nThe PPM information-cells are decoded by means of a head sensitive to the field or rate of change of the flux emanating from the polarized regions at the transition zones. The signal produced by most magnetic reproduce heads is a voltage pulse corresponding to the location of a transition zone from one magnetically set region to another. The detection circuitry to which the head is attached must determine the location of one peak relative to another. This is most often done by differentiating the pulse and comparing the voltage to zero. The differentiated waveform will cross zero at a time corresponding to the peak in the original waveform, at which point the comparator will change state. This works well in the region where the pulse has energy. However, outside this region, the input to the comparator hovers around zero, crossing zero as a result of noise, producing transitions at the output. To remedy this, the output of the comparator is gated by a circuit which only allows the comparator state to change when the input waveform exceeds a threshold, usually 30%-45% of the peak amplitude. This scheme works well in eliminating the unwanted transitions. A side effect is that, if the peak amplitude of the input signal fails to reach the threshold, no transition will occur and the data will be missed. Hard disk drives deal with this problem by marking locations on the disk where this occurs and not using the areas. Tape drives use frequent re-synching, read-while-write, re-reading, and powerful error correction code (ECC) to deal with the problem. Simple, inexpensive systems for use in cameras cannot use such complex and expensive solutions, and resort is therefor made to the self clocking PPM coding of Wash or Chi.\nAs shown in FIG. 1B, the data and clock transitions are read out as positive and negative going spike signals at the transitions. Turning to the read circuit 10 of FIG. 2, reproduced from the '139 patent, a read/write magnetic head 12 provides the read signals of FIG. 1B on lines 14. A pre-amplifier 16 amplifies the output signal from the magnetic head 12 and applies the amplified signal on lines 26 to a filter 18 for removing an unnecessary component from the amplified read signals. A post-amplifier 20 amplifies the filtered read signals and applies the filtered and amplified, positive and negative, signals to a peak detector 24 via line 22. In detector 24, the signals are applied to a positive peak detector (NPD) 30 and a negative peak detector (PPD) 32 for detecting the negative and positive pulses having amplitudes exceeding respective threshold levels to ensure that noise signals are not decoded as data or clock signals. Uniform amplitude, polarity, and pulse width, clock and data, signal pulse trains are thereby generated on lines 38 and 34 as depicted in FIG. 1C. A delay circuit 40 is connected to the NPD 32 via line 38, resulting in a delayed clock pulse train on lines 46 and 48 and also appearing in FIG. 1C.\nThe delayed clock signal is applied to the set terminal of the flip-flop 36 and is also applied to clear the up/down counter 50. The data signal is applied to the clear terminal of the flip-flop 36, so that the flip-flop 36 is set by delayed clock signals and cleared by data signals to provide the square wave signal shown in FIG. 1D at the Q output. The Q output is set high by a delayed clock signal and low by the succeeding data signal.\nDuring the high state, the up/down counter 50 counts system clock 56 pulses applied to its CLK input, incrementing the count when the Q output is high and decrementing the count when the Q output is low. The most significant bit (MSB) of the output count of the up/down counter 50 and the CLOCK output from the NPD 32 are supplied to a computer 42. The particular manner of determining the data bit in the information-cell is further explained in the '139 patent.\nBecause of the pulse position modulation, variation in velocity of the filmstrip during recording or readout (within reasonable bounds in which a camera would be specified to operate) does not effect the accuracy of the data, as described above. However, the compliance of the magnetic record/reproduce head with the low density MOF layer during recording or reproducing the recorded data itself may be faulty for other reasons leading to a failure to either read or write data in one or more of the information-cells.\nWhen compliance with the magnetic head is lost due to a dirt speck on the magnetic medium or curl of the medium or jitter in the head suspension or if the magnetic recording medium has a defect such as a change in the density of the magnetic powder or if electrical noise is introduced in a playback system from outside the system, false data generally referred to as a \"drop-in\" and/or the omission of data generally referred to as a \"drop-out\" can occur in the reproduced data stream. As described above with respect to FIG. 2, the conventional data read out circuit compares the absolute positive and negative amplitudes of the reproduced analog waveforms to threshold levels in the PPD and NPD circuits to convert, among the above-mentioned peaks, only the peaks having amplitudes higher than a predetermined threshold level into digital signals, determining that they are data bits 1 or 0 depending on their position in the case of the PPM encoded signals. Thus, if the amplitudes of the pulse signals fade below the thresholds, both the clock pulses separating successive information-cells and the data pulses signifying the binary data content can be lost.\nThe information-cells are recorded in a pattern of multi-bit bytes described in detail in the '419 patent, particularly in regard to FIGS. 6-9 thereof. A simplification of that format appears as follows in table I:\n______________________________________ Original Data ______________________________________ P.sub.1 d.sub.11 d.sub.12 d.sub.13 d.sub.14 d.sub.15 d.sub.16 d.sub.17 P.sub.2 d.sub.21 d.sub.22 d.sub.23 d.sub.24 d.sub.25 d.sub.26 d.sub.27 P.sub.3 d.sub.31 d.sub.32 d.sub.33 d.sub.34 d.sub.35 d.sub.36 d.sub.37 P P P P P P P P ______________________________________\nA block of three data bytes of eight information cells is shown. The first bit of each byte is a parity bit (P.sub.i), and the remaining seven bits are data bits (d.sub.i1 -d.sub.i7). Each block of bytes recorded serially in a single track in the MOF layer is separated by a parity byte, called a Longitudinal Redundancy Check (LRC). Each vertical bit value of the LRC is calculated from the values of the data bits in its respective column. The data bits of successive information-cells in each block are thereby protected with a 2-dimensional, horizontal and vertical, parity check in the data blocks. This scheme is able to correct an odd number of bits in error in any one byte, provided a bit is present at every position and the LRC is present in its proper position with the longitudinal parity bits lined up with their respective columns.\nThe PPM encoding method described above incorporates the clock information in with the data so that a drop-out typically results in both loss of data content and loss of the clock signal. When a drop-out occurs, no clock is available to transmit data from the decoder to the memory, so that bits are completely missing in the recovered data block. This results in loss of byte synchronization and renders the parity checks useless in recovering the uncorrupted data. Shown below in table II is the situation after recovery when bits are lost:\n______________________________________ Corrupted Data ______________________________________ P.sub.1 d.sub.11 d.sub.12 d.sub.13 d.sub.14 d.sub.15 d.sub.16 d.sub.17 P.sub.2 d.sub.21 d.sub.22 d.sub.23 d.sub.24 P.sub.3 d.sub.31 d.sub.32 d.sub.33 d.sub.34 d.sub.35 d.sub.36 d.sub.37 P P P P P P P P X X X ______________________________________\nA drop-out of three information-cells and the respective data bits d.sub.25, d.sub.26, d.sub.27, occurred in the second row, causing the data bytes and the LRC parity bits to shift in the recovered block. No correction is possible because of this, and the data read out may either be thwarted or the data may be misinterpreted. It should be understood that in practice, the problem is magnified by the number of bytes in the data set and the location of the drop-out in the series of bytes."}
{"text": "This invention relates to a nozzle for a turbojet engine and, more particularly, to such a nozzle that is convergent-divergent, has a variable throat, and is air cooled (i.e., the hot gas stream surface of the nozzle is cooled by air).\nThe prior art method of cooling the convergent and divergent hot gas stream surface of a nozzle comprised cooling the convergent portion of the nozzle by using a liner, which was suported above the convergent portion of the nozzle, thus forming a passage for the cooling air (i.e., the fan duct air). The convergent cooling liner(s) extended from the end of the augmentor liner to the nozzle throat. The cooling air was then \"dumped\" onto the divergent section, to \"film cool\" that section.\nThere are two principal significant problems with regard to this method, and the structures used to perfom the method. Firstly, the convergent liner(s) warp severely, due to augmentor temperature gradient, thereby creating undesired leakage paths for the cooling air. Secondly, the cooling air boundary layer has difficulty turning around (i.e., going over) the sharp throat radius, thereby resulting in severe warping of the divergent section surface(s), preventing the divergent surface(s) from effectively forming an interface seal to hot gas stream leakage, and thereby creating performance losses and very short life for the convergent section liner(s) and the constituent portions of the divergent surface(s).\nOur invention obviates these well known, and long-standing, problems of the prior art; and, therefore, our invention significantly advances the state-of-the-art."}
{"text": "(1) Field of the Invention\nThe present invention relates to encapsulated foods filled with a desired edible liquid and the process for preparing drinks which include such encapsulated foods. Furthermore, this invention relates to a process for preparing encapsulated foods to make, encapsulated dressings, encapsulated alcohol drinks, encapsulated coffee or red tea etc. by using dressings, fruit juice, alcoholic drinks and other liquids, popular foods, liquid sweetening materials, etc. as edible liquids and the process for preparing drinks which include these encapsulated foods.\nGenerally, popular foods such as alcohol, coffee, red tea, cocoa, etc. as well as fruit juice are used as drinks because of their liquid form. Moreover, acidic condiments, i.e. dressings and condiment vinegar etc., being normally liquid, are poured over a salad, for example, by using a spoon. However, the poured condiments either permeate into the salad itself or are dispersed and absorbed into any food surrounding it. Therefore, it is impossible to visually confirm the quantity that has been used or to know for certain how much of it enhanced the flavor of the salad. Although powdered condiments and granular ones may resolve these problem to a certain extent, the use of powdered ones is somewhat restricted due to their quality or taste.\nThe present invention intends to resolve these problem. More specifically, it provides a new process for preparing edible dressings, edible alcoholic drinks, edible coffee or tea, etc. which have the taste and quality of liquid dressings, fruit juices, alcoholic drinks, popular liquid foods, liquid sweetening materials, etc., and besides have more visual character and are not too soft.\n(2) Description of Prior Preparation\nPreviously, a process for preparing capsules using solutions of alginic acid salt and calcium salt is conducted by dropping solutions of alginic acid salt and calcium into another solution to form a membrane of calcium alginate, thereby forming the capsule. However, it is impossible to obtain suitable capsules when this process is applied as it is to dressings, alcoholic drinks or to any preferred drinks.\nFor example, when an acidic condiment per se is subjected to encapsulation, the equilibrium of the alginic acid ion and sodium ion is lost by the acidifying the sodium alginate solution to agglomerate and generate alginic acid separately. Also, the calcium ions for forming the membrane are fewer because they react with acid so as to extract calcium salt by acidifying the calcium salt solution; therefore, in any way, it is difficult to produce good capsules.\nIt may be considered that the process is to first form the capsules and then fill them with an acidic condiment liquid. However, in this case, the capsules consisting of calcium alginate are influenced by the acid due to a time lapse, thus causing various problems to arise such as how to produce capsules with strong membranes which will neither collapse nor change in quality under preservation when filling them with the acidic condiment liquid and how to fill the acidic condiment liquid into the inner spaces of the capsules already formed.\nWhen fruit juice per se is encapsulated as a core material, various problems arise such as that fruit juice reacts with alginic acid to increase its viscosity and alginic acid generates and agglomerates separately since fruit juice is normally an organic acid liquid of pH 3-4. Furthermore, it becomes difficult to conduct a strong bridge reaction in view of the competing COOH group of alginic acid molecules and the COOH group of fruit juice molecules because fruit juice reacts with calcium salt to decrease the number of calcium ions and further generates alginic acid calcium, producing an inferior taste so that, being not edible, this causes a problem that is difficult to solve.\nConcerning alcoholic drinks to be encapsulated as core material, here also various problems arise making it difficult to form the desired capsules despite the addition of products such as xanthene gum, etc.; besides, drops of alcoholic drinks collapse on the surface of alginic acid salt liquid when they are dropped into it. The reason for such a phenomena is due to the fact that alcoholic drinks per se are low in molecular compounds and have a low viscosity. In any case, there is also the problem of solving the difficulties in forming capsules by such a reaction.\nThe same problem exists in the case of other liquids and liquid sweeters. Namely, when these substances are encapsulated as core material those in which the mixture of alginic acid salt with liquids or liquid sweetening material is dropped into calcium salt water solution may be encapsulated well; however, this manner of preparation is not desirable for taste because the resulting capsules become hard. Furthermore, when the mixture of calcium salt water solution and liquid sweetening material is dropped into an alginic acid water solution, an unsuitable taste is generated by the remaining calcium constituent in the sweet materials. In any case, it is difficult to prepare suitable capsules. Such problems also occur with encapsulated foods including cut solid foods containing the core liquid."}
{"text": "1. Field of the Invention\nThis invention relates to an optical exposure apparatus for projecting a reduced pattern of an original mask pattern precisely onto a desired position of a workpiece to be exposed, and more particularly to an improvement of an alignment system of such an apparatus for precise alignment of the original mask pattern and the workpiece.\n2. Description of the Prior Art\nSemiconductor devices, such as transistors or integrated circuits, are usually fabricated by a series of diffusion steps. Each diffusion step involves applying a coat of photosensitive polymer, known as photo-resist, over a silicon dioxide layer on the surface of the semiconductor substrate. Then, the photo-resist is exposed to light through a mask having a predetermined pattern by means of an optical exposure apparatus. When the photo-resist is photographically developed, selected areas of the photo-resist are removed to expose the underlying silicon dioxide. The exposed silicon dioxide is then removed by an etching fluid which does not attack the photo-resist to expose the underlying semiconductor material. The photo-resist is then stripped from the silicon dioxide and impurities are diffused into the areas of the semiconductor material exposed by the openings in the silicon dioxide layer. A new silicon dioxide layer is either grown over the exposed portion of the semiconductor material during the diffusion process, or is subsequently deposited, and the above-mentioned steps are repeated as often as required for the formation of the desired circuitry.\nAs the above-mentioned optical exposure apparatus, a so called \"step and repeat camera\" is usually used. A schematic block diagram of the step and repeat camera of the conventional type is shown in FIG. 1.\nReferring now to FIG. 1, the camera is principally composed of a light source 6, a condenser lens 7, a mask 1, a reduction lens 8, a moving table 4, and a table position control circuit 3. In actuality, the moving table 4 can be moved in the X and Y directions, but, for simplifcation, the table movement in the Y direction is not shown in FIG. 1.\nIn operating the camera, firstly, the mask 1, which has an original pattern 1a formed by an opaque material, is set on a stationary holder H.. And, a semiconductor wafer 2 is placed on the table 4. Then, the table position is adjusted so that the desired exposure position of the semiconductor wafer 2 is precisely located at the projection position of the reduced pattern of the original pattern. That is, the table position control circuit 3 controls the table position in response to a control signal S1 supplied from a control signal generator 9, so that a proper alignment is obtained. When this operation for controlling the table position has been finished, a signal S2 is delivered to an exposure control circuit 5, and then the light source 6 begins to operate in response to an exposure control signal S3 supplied from the circuit 5.\nA dispersed light beam l.sub.1 from the light source 6 is formed into a parallel light beam l.sub.2 by the condenser lens 7, and then the mask 1 is irradiated by this parrallel light beam. A light beam l.sub.3 passing through the mask 1 is projected on the wafer 2 through the reduction lens 8, so that the photo-resist layer disposed on the wafer 2 is exposed by the reduced pattern of the original pattern.\nIn the above-mentioned conventional optical exposure apparatus, the accuracy of the alignment of the exposure position and projection position is decided approximately according to the accuracy of positioning the table 4. Therefore, in order to increase the accuracy of the fabrication of the semiconductor devices, it is necessary to provide a X-Y moving table having a very high accuracy. As a result, the moving table becomes very expensive and heavy in weight. Additionally, the moving speed of the table is restricted to small a value because of its heavy weight. In a result, the time period required for one time of exposure becomes long, and hence the efficiency of the fabrication decreses. Furtheremore, it is impossible to ignore the decrease of the accuracy of the table positioning under the influence of the changes of the circumferential conditions, such as temperature, during the exposure step."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an image pick-up tube and particularly to an image pick-up tube which can easily be manufactured and of which the characteristic such as aberration or the like can be improved.\n2. Description of the Prior Art\nFIGS. 1 and 2 respectively show examples of the prior art image pick-up tube. In FIGS. 1 and 2, reference numeral 1 designate an electron gun and 2 a glass envelope. In FIG. 1, reference characters G3, G4 and G5 respectively designate grid electrodes from which leads 3, 4 and 5 are respectively led out to apply predetermined voltages to the grid electrodes G3, G4 and G5. In FIG. 2, reference characters G3 and G4 respectively designate grid electrodes from which leads 6 and 7 are respectively led out to apply predetermined voltages thereto. In the image pick-up tube shown in FIG. 2, the electrode G4 is formed by, for example, depositing metal on the inner surface of the glass envelope 2.\nSince the image pick-up tubes having the configuration shown in FIGS. 1 and 2 each employ the cylindrical electrodes, such prior art image pick-up tubes require a mechanical construction (not shown) to support the cylindrical electrodes and leads must be connected (welded) to the electrodes thereof, thus the manufacturing of the image pick-up tube being made difficult. In such prior art image pick-up tubes, electrostatic lenses are formed between the electrodes. In this case, since the image pick-up tubes having the configuration shown in FIGS. 1 and 2 employ the cylindrical electrodes apart from the inner surface of the envelope 2, the aperture of the electrostatic lens thus formed is smaller than the inner diameter of the glass envelope 2 so that distortion becomes serious and thus the characteristics such as aberration or the like are deteriorated."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to a container for powder cosmetic products. More particularly, it relates to a cosmetics container that provides for a leakproof storage and ease of application of a powder cosmetic product. The container comprises a freely rotatable telescopic applicator that is integral with the cap and said container with a double walled cavity is having a non-circular aperture inside, which is plugged by the rod of the applicator element. Further, the cap of the container is capable of securely fitting onto the neck of the container by way of threads in the cap, thereby providing a leaktight storage and convenient use thereof.\n2. Description of the Related Art\nCosmetic containers for powder cosmetics are known which generally include a container for holding the product having a wiper element in which sits the applicator in the form of a brush formed by a handle carrying a tuft of bristles. The container is provided with an additional cap to be screwed on to the container. Also, the applicator is of the size of the container which is generally small and is therefore inconvenient in application. Such a container has inconvenience in using as it consumes a lot of time by the time the product is ready to be applied. The shortcomings include time consumption in shaking the container, separating the cap from the case and then withdrawing the brush, possibility of misplacing the cap or the applicator if the package is left in disassembled state after completion of its use. Moreover, when the device is shaken, it is possible for powder to escape from the inside of the container to be deposited around the brush handle, inside the cap. The user's fingers are then soiled with product when he/she grasps the handle.\nThere have been proposed improvements over the aforementioned containers to overcome the shortcomings of the conventional case for powder cosmetics. For example United States Patent Application Publication No. 2005/0092343 discloses a cosmetics case for storing color cosmetics in powder state. The disclosed cosmetics case includes a cosmetic aid covered with powdery cosmetics and capable of being fixed to or released from a case body to be freely used by turning a cover of the cosmetics case to the right and left directions within a range of 90 degree angles. Although the case can eliminate the inconvenience caused due to loss of the cover in conventional products, the problem of the user's fingers getting soiled with the product still remains. Moreover, if the user forgets to rotate the cover to lock the container, the problem of leakage still exists.\nUnited States Patent Application Publication No. 2006/0239753 discloses a device for packaging and applying a product that includes a container holding the product to be applied in loose powder form, the container having an opening. The device also includes an applicator including a rod, an applicator element at a first end of the rod, and a grasping portion connected to a second end of the rod. The applicator is configured to be attached by snapping onto the container so as to accommodate the applicator element in the container. However, as the cap is not screwed onto the container, the chances of leakage and the user's hands getting soiled during use still exist.\nAlso known are various containers wherein telescopic applicators are fixed to the caps of the containers in a spring-biased manner, however, the length of the rod of such telescopic applicators extends to bottom of the container such that the applicator element touches the bottom of the container in order to securely close the container. The drawback associated with such an arrangement is that there is a possibility of the applicator getting damaged with usage as the applicator is always under strain when in a closed position. For example, U.S. Pat. No. 6,883,990 discloses a cosmetics container comprising a hollow container and a cover unit. The cover unit covered on the upper end of the hollow container is composed of a connecting cover, a brush, a revolving base and a revolving cover, with the brush positioned in the connecting cover. While using, the revolving cover is to be turned around to move forward together with the revolving base and actuate the brush to move forward together and extend out of the cover unit for use.\nAlthough many of these prior art containers are relevant with respect to the present invention, most of these containers disclose that the applicator tip sits into the wiper element and thus there is a possibility of the applicator being damaged in further uses as the applicator is always under strain. Moreover, the aforementioned products have been unable to provide the user with a container that allows an easy and convenient storing and applying of a powder cosmetic product. Also, when the applicator is integrally attached to the cap and the applicator plugs into the wiper element, it is necessary that the applicator and the wiper be circular in cross section so as to allow the cap to be screwed onto the container. Otherwise, the cap has to snap fitted on the container which again has the shortcoming of not being a leak proof closing for a container storing powder cosmetics.\nTherefore, there is a need in the art for a container for powder cosmetics container that overcomes the shortcomings of the known containers. It is also desirable that the container provides ease of storage and application of the product."}
{"text": "The present invention relates to xerography or printing systems, in particularly, those systems having a color print output. More specifically, the subject invention relates to a user interface for selectively adjusting hue of a color print. However, it will be appreciated to those skilled in the art that the invention could be readily adapted for other adjustable features of a print or display output as, for example, chroma, luminescence, or where similar user interfaces are employed to implement operator preferences.\nThe printing or display of image information via a printer, copier, or personal computer can be controlled by an operator through a variety of interfaces. Common interfaces comprise a number of systems and devices such as a combination touch pad and display screen, a keypad or a mouse control often accessing a tool bar for icon display of system features and options, all having some capacity for controlling a printer or display. To facilitate operator ease of use of the interface controls, such systems are known to include a representative illustration identifying a consequence of a control command. Such illustrations typically comprise an icon which will represent a state or change of state in response to a user instruction. A common example is a representation of paper size in a copying system as the user adjusts printing output from letter, to A4 to legal or other sizes. A problem with such representative icons is that they fail to communicate an accurate identification of consequences of subtle appearance changes in the output such as may occur in color printing from minor adjustments of the input controls. In particular, an icon representing a system state may fail to adequately communicate an accurate indication of the system output. On the other hand, where a number of controls are being adjusted, an icon identifying relative adjustment away from a preexisting default state, can merely identify a tendency away from the default state, but also cannot accurately and adequately represent the consequence on the output. For example, adjustments in hue control such as hue shift are usually available to a user to satisfy a particular user\"\"s preferences. The phrases xe2x80x9chue shiftxe2x80x9d or xe2x80x9chue controlxe2x80x9d are used to describe the changes in the attributes of the perceived color by an operator throughout a full color spectrum. Such a spectrum is directly related to different light wave lengths as perceived by a human user and, for example, will comprise a continuous color spectrum comprising red, yellow, green, blue, purple and magenta. See Billmeyer and Saltzman, Principles of Color Technology, Second Edition, Chapter 2, xe2x80x9cDescribing Color xe2x80x9dJohn Wiley and Sons, Inc., 1981.\nFIG. 1 shows a prior art implementation of a control interface for a color copier for shifting color hue. The interface comprises a screen on an LCD display also including touch screen controls. In this implementation, concentric color spectrum rings are used to suggest an effect of a user adjustment. More particularly, the inner ring 10 defines a reference of the red, yellow, green, cyan, blue and magenta spectrum as an indicator of how the colors of the original input image should be perceived by the copier system.\nThe outer ring 12 comprises an identical spectrum which at a default state is exactly aligned with the inner ring to indicate that the output print will have an identical system hue as the input hue. However, a user upon checking the print output, may determine that the output hue is unacceptable or improvable. As a practical matter, most users will endeavor to shift color hue in an effort to improve skin tone of the output print image, but present systems provide no representative image which is more accurately indicative of how the shift will affect the overall output, and the user will often have to step through a number of copy operations until, by trial and error, an acceptable print output is obtained. Without such trial and error operations, it is very difficult for an operator to envision changes to a color hue as they are implemented throughout the entire spectrum of output print colors merely as a consequence of trying to improve a skin tone so that such effort can be a tedious problem to overcome.\nMore specifically and with reference to a copying operation, this feature allows a shift in all of the colors towards the adjacent hue. Another example of where this can be useful is when the document being copied does not have quite the correct colors for the application. By touching one of the two arrow buttons (14, 16) one can cause the overall copy colors to be shifted.\nThe inner color ring 10 represents the colors of the document being copied and rotates to indicate the direction in which the colors shift. The outer ring 12 represents the colors that appear on the copy (not shown). For example, if the right arrow button 14 is selected, the red colors R shift towards the yellow hue Y, blue colors B shift towards the magenta hue M, and green colors G, shift towards the cyan hue C. All colors in between these colors also shift in the same direction.\nIf the left arrow button 14 is selected, the red colors R shift towards the magenta hue M, blue colors B shift towards the cyan hue C, and green colors G shift towards the yellow hue Y. Again, all colors in between these colors also shift in the same direction.\nThe interface system of FIG. 1, although suggesting a shift in color hue, fails to apprise a user of an accurate consequence of the shift for an adequate determination of how much or how little shift is appropriate. Trial and error must be used.\nThe present invention contemplates a new and improved hue shift control system which overcomes the problems referenced above and others to provide a new color controller which is simple in design, readily understandable by a user, easy to operate, and which provides improved communication of hue shift commands with reference to a dynamic image representation of a print output.\nIn accordance with the present invention, there is provided a control system for a color printer or copier including a user interface for selectively adjusting a print output. The system comprises an illustrative representation of an adjustable color spectrum suggesting to the user a color range for the output. The user interface also has a selective control including an indicator of the user input command and for representing an effect of the user command on the output. A dynamic image representation suggesting the effect on the print output concurrently communicates an adjustment in the appearance of the print output in accordance with the implementation of the user commands. The image representation thus provides a good indicator to the user of the effect of a change in hue command across the whole entirety of the color spectrum so that the user can more quickly and efficiently identify the likely effects of such control commands on the print output.\nIn accordance with another aspect of the present invention, the selective control comprises a slide bar suggestive of the color hue spectrum. A touch-sensitive mechanism overlay allows the user to implement commands so that the slide bar moves relative to the representation of the color spectrum to suggest to the user the consequences in hue shift as a result of the input control command.\nIn accordance with a further aspect of the present invention, the dynamic image representation includes a skin tone portion displayed in an environment substantially including the rest of the color spectrum. The consequence of hue shift commands can thus be perceived by the user not only with reference to the skin tone of the image representation but of the effects on other portions of the image representation.\nIn accordance with another aspect of the present invention, both the indicator and the dynamic image representation concurrently communicate consequences of the user commands to jointly indicate to the user the probable consequences of the commands on the print output.\nIn accordance with yet another aspect of the present invention, the interface includes a selective default setting for a relative setting of the hue control.\nOne benefit obtained by use of the present invention is a control system wherein the user interface provides a dynamic image representation as part of a display or control panel wherein a direct indication is displayed of the effect of a change or shift in the feature of the output resulting from the user command. The image is represented in concert with a reference representation of a spectrum of printed attributes for accurately communicating an identified change in a basic spectrum relative to the associated image representation.\nAnother benefit obtained by use of the present invention is a control system which provides improved communication of the effects in hue shifts. When a user shifts overall copier or print color to an adjacent hue, the effect of the color shift on the entire spectrum is appreciable along with the likely consequence of the change in appearance of a print output.\nA more limited benefit of the subject invention is that the controller has a default position for a predetermined standard position for the controller, which standard position also is easily understood and represented by the user interface and image representation.\nOther benefits and advantages for the subject new user interface and printer controller will become apparent to those skilled in the art upon a reading and understanding of the specification."}
{"text": "The subject matter discussed in the background section should not be assumed to be prior art merely as a result of its mention in the background section. Similarly, a problem mentioned in the background section or associated with the subject matter of the background section should not be assumed to have been previously recognized in the prior art. The subject matter in the background section merely represents different approaches, which in and of themselves may also correspond to implementations of the claimed technology.\nThe Naive Bayes Classifier (NBC) is a stalwart of the machine-learning community, often the first algorithm tried in a new arena, but also one with the inherent weakness of its hallmark assumption—that features are independent. Bayesian networks relax this assumption by encoding feature dependence in the structure of the network. This can work well in classification applications in which there is substantial dependence among certain sets of features, and such dependence is either known or is learnable from a sufficiently large training set. (In the latter case, one may use any of a number of structure-learning algorithms; they don't always work well because the structure-learning problem is very difficult.) Undirected graphical models are a different way to relax this assumption. Like Bayesian networks, they also require domain knowledge or a structure-learning procedure. Domain knowledge can be hand coded into the structure, but this only works when domain knowledge is available and even then hand coding is usually very laborious. Alternatively, structure-learning algorithms require extensive training data and present a computationally complex problem, in addition to the concern of over-fitting a model using a limited data set, thereby creating a model that predicts structures that are less likely to work well on data not seen during training Undirected graphical models may also fail to work well when either the domain knowledge is wrong or when the instance of the structure-learning problem is either intrinsically hard or there is insufficient data to train it.\nThe technology disclosed relates to machine learning (ML) systems and methods for determining feature dependencies—methods which occupy a space in between NBC and Bayesian networks, while maintaining the basic framework of Bayesian Classification.\nBig data systems now analyze large data sets in interesting ways. However, many times systems that implement big data approaches are heavily dependent on the expertise of the engineer who has considered the data set and its expected structure. The larger the number of features of a data set, sometimes called fields or attributes of a record, the more possibilities there are for analyzing combinations of features and feature values.\nAccordingly, an opportunity arises to automatically analyze large data sets quickly and effectively. Examples of uses for the disclosed classification systems and methods include identifying fraudulent registrations, identifying purchase likelihood for contacts, and identifying feature dependencies to enhance an existing NBC application. The disclosed classification technology consistently and significantly outperforms NBC."}
{"text": "1. Technical Field of the Invention\nThe present invention relates generally to communications networks. More particularly, the present invention relates to a system and method for optimizing the deployment of network elements over a period of time in a network, e.g., a Fiber Optic Network, that routes bandwidth demands having multiple channel rate requirements and variable time points. In addition, the present patent application provides a scheme for routing demands in a Fiber Optic Network using Time Slot Assignment (TSA) technology.\n2. Description of Related Art\nInstalling, maintaining, and upgrading a communications network is very costly. Huge expenditures are involved in deploying suitable network equipment at predetermined locations and establishing transmission paths or conduits therebetween via an appropriate physical medium (or, media). Because of these cost considerations, network operators have to be circumspect about when and where to put in a new network or expand an existing one.\nFurthermore, information transport—the primary function of a communications network—needs to be efficiently performed in a network in order to optimize its size so that costs associated with unnecessary expansion of the network, sub-optimal deployment/upgrading of the equipment, etc., for example, are avoided. As is well known, efficient network usage generally implies using available channel capacity or capacities efficiently, in addition to employing techniques that achieve cost-effective routing of information via available network equipment.\nIt should be readily apparent that there is a need for planning tools that enable network operators and owners to schedule equipment deployment intelligently, especially in light of the aforementioned considerations. In addition, such tools have lately become even more essential because of the ever-increasing need for the deployment of high-capacity networks (thus involving more sophisticated and expensive equipment) capable of transporting a wide variety of information—voice, data, video, multimedia, and the like—at phenomenal transport rates.\nConventional solutions in this regard typically employ mathematical modeling or simulation techniques coupled with optimization procedures to arrive at estimates for placement of network equipment that routes information as efficiently as possible. Although such methodologies represent a significant advancement in the field of communications network modeling, there exist several shortcomings and deficiencies in the state of the art.\nFirst, the existing methods treat demand, a quantified volume of bandwidth requested to transfer information over a network path, as a time-independent parameter, thereby compressing all demands—current and projected—to be serviced by the network to a single point in time. In other words, all network equipment required to satisfy both current and projected demands is treated as operational at a single instance. Those skilled in the art should readily recognize that while such a technique may yield a “good” first approximation, it is nevertheless unsatisfactory for accurate planning purposes where new network portions (e.g., rings) are built in a phased manner across the life of a deployment plan, typically stretching over several quarters or years.\nFurther, as a by-product of treating demands as time-invariant entities, resultant mathematical formulations become formidable because, typically, several hundreds of thousands of demand quantities—including demand forecasts—need to be optimized (that is, demands to be optimally routed in a network) over a deployment plan. Computation loads therefore become enormous, leading to critical time delays in obtaining results which often tend to be unstable because of the unwieldy modeling efforts.\nIn addition, the existing solutions typically consider only a single type of channel bandwidth for the demand quantities that need to be optimized. Moreover, the channel rate thus considered is oftentimes a lower rate, thereby necessitating decomposition of demands of higher channel bandwidth rates into equivalent demand units of the lower rate used. However, no controls are implemented to ensure that these equivalent lower rate demands are routed together on the same network paths or to the same intended destinations. Clearly, such routing is unacceptable and is only a poor approximation of the actual routing loads in the network.\nYet another drawback in the current network-planning methodologies is where the underlying modeling apparatus does not accurately reflect today's network transport technology. For example, where Fiber Optic Network rings are implemented, current solutions yield results which are not compatible with the transport technology that is widely deployed."}
{"text": "It is known that fuel level and composition information can be collected from a combination of capacitors located inside a fuel tank. The capacitance of a capacitor immersed in a fluid is proportional to the dielectric constant of the fluid times the capacitor's free air capacitance according to the following equation: EQU C.sub.f .apprxeq..epsilon..times.C.sub.a ( 1)\nwhere:\nC.sub.f =capacitance (pF) of capacitor C immersed in fluid PA1 .epsilon.=dielectric constant of the fluid in the tank PA1 C.sub.a =capacitance (pF) of capacitor C in free air PA1 C.sub.f =capacitance (pF) of capacitor C immersed in fluid PA1 .epsilon.=dielectric constant of the fluid in the tank PA1 C.sub.a =capacitance (pF) of capacitor C in free air PA1 L=total length (inches) of capacitor C PA1 H=height (inches) of the fluid in the tank\nA method for sensing the level of a fluid in a fuel tank can be derived using equation (1) and the measured value of the capacitance developed on a partially immersed capacitor in a fluid of a known dielectric value. For example, U.S. Pat. No. 5,051,921, issued Sep. 24, 1991 to Paglione, describes such a method which uses interdigitated capacitors immersed in fluid contained in a fuel tank to sense the level of the fluid.\nUsing capacitive fuel gauges to detect the level of a fluid other than pure gasoline however, creates a problem since most capacitive fuel gauges are calibrated according to the dielectric properties of pure gasoline. An output signal from a capacitive fuel gauge immersed in a mixture of fuel ingredients, e.g. a mixture of gasoline and methanol or ethanol, registers a false reading if the gauge was calibrated to sense the level of pure gasoline since the gasoline mixture has a dielectric constant different from that of pure gasoline.\nThe method of the '921 patent addresses the capacitive fuel gauge calibration problem by first recognizing that the level of any fluid in a tank can be calculated according to the following equation: EQU C.sub.f .apprxeq.C.sub.a +(C.sub.a /L) [H(.epsilon.-1)] (2)\nwhere:\nAs described in the '921 patent, C.sub.f, C.sub.a, L are known quantities. .epsilon. is determined using equation (1) and measured capacitance values obtained from a capacitor immersed in the fuel. The value of each of these parameters is substituted into equation (2) to determine H as follows: ##EQU1## The fuel level detection scheme disclosed in the '921 patent utilizes two capacitors connected to a monostable multivibrator, also known in the art as a \"one-shot,\" and an oscillator or clock. The monostable multivibrator modifies the width or duration of output pulses from the oscillator or clock in proportion to the capacitance of each capacitor. The output of the first capacitor, directly proportional to the dielectric constant or composition of the fluid in the tank as indicated by equation (1), is called the composition signal. The output of the second capacitor, directly proportional to the level of the fluid in the fuel tank as indicated in equation (2), is called the level signal. An analog processor multiplies the level and composition signals together to obtain a fuel level signal which incorporates the correct dielectric constant of the fuel mixture. Thus, a signal is produced which always represents the correct fluid level regardless of fuel composition.\nAlthough the method described in the '921 patent solves the problem of obtaining an accurate fluid level reading for mixtures other than pure gasoline from fuel gauges calibrated according to the dielectric constants of pure gasoline, it employs expensive components such as multi-vibrators, an external comparator, multiple low pass filters and an analog math co-processors which make the '921 method too expensive for use in production vehicles.\nIn addition, the method disclosed in the '921 patent proves disadvantageous due to the location of the composition capacitor at the bottom of the fuel tank. Water and contaminating particles inevitably find their way to the bottom of the fuel tank. The accumulation of pollutants tends to corrupt the capacitor changing the amount of charge that the capacitor stores which results in inaccurate readings.\nFurthermore, the '921 patent does not provide for free-air calibration of the fuel gauge before the unit is submersed in the fluid making the '921 device very sensitive to the manufacturing tolerances of the various component parts.\nMoreover, the multiple sensors used to sense level and composition data have traditionally been connected to their individual displays using individual wires. The individual wires have either been banded together into a wiring harness that exits the fuel tank through a relatively large opening or each wire has been routed out of fuel tank individually through one of many openings in the fuel tank. Using a scheme that integrates the various output data signals onto a single wire for transmission to a central vehicle computer via a multiplex bus reduces the number of openings in the fuel tank through which harmful emissions can emanate.\nFinally, integration of the various functions into one fuel gauge assembly eliminates redundant subsystem components, such as multiple controllers, resulting in a significant reduction in the cost of the fuel gauge assembly and in its complexity which, in turn, facilitates the manufacturing process."}
{"text": "Direct conversion receivers (DCRs) are prevalent in mobile communication systems due to their simplicity, low cost and for the reason that much of the signal processing can be achieved in the digital domain. A drawback to the elimination of an intermediate frequency stage is an increased susceptibility to second-order intermodulation effects. For example, intermodulation products generated from strong interfering signals can be downconverted to baseband with relatively weaker signals of interest, thereby desensitizing the DCR to the desired signal. Such interfering signals, referred to herein as blocking signals, are ubiquitous; they may originate from communication signals in adjacent channels and/or from transmission sources that, even when far removed in frequency from the channel of interest, transmit at a power level sufficient to cause significant interference problems in a DCR, including bit error rate (BER) and/or signal-to-noise ratio (SNR) degradation.\nIn view of the susceptibility to distortion by second-order effects, design specification of DCRs typically include high input-related second order intercept point (IIP2) criteria. IIP2 is the theoretical input signal level at which second-order intermodulation products are equal in power to that of a desired signal. Thus, if the IIP2 can be made higher, the power of an interfering signal must reach a correspondingly higher level to have an equivalent detrimental effect on the DCR. It is thus clearly desirable to establish IIP2 at as high a level as possible.\nSecond-order intermodulation distortion is an effect of non-linear behavior of components in a DCR and downconverting mixers are often mostly responsible for limits on the level to which IIP2 can be established. Compensation techniques can be applied to increase the IIP2, but changes in temperature and/or frequency can reduce the IIP2 to effectively undo applied compensation measures. Indeed, an IIP2 established at 100 dBm can be reduced to 50 dBm in response to a change in frequency or temperature.\nGiven the potential improvements in DCR performance that can be expected upon overcoming the foregoing limitations in the art, substantial resources are increasingly devoted toward developing interference-tolerant designs in direct conversion receivers."}
{"text": "Retaining walls are used widely to construct soil embankments, sound barriers, and highway dividers, as well as for other similar and dissimilar purposes.\nA retaining wall of a type in widespread use is made from concrete modules stacked in multiple courses. The modules are spaced laterally from one another in each course so that modules of each intermediate course are straddled by modules of the next course above and by modules of the next course below.\nSince the modules are stacked so as to provide spaces therewithin and therebetween for plantable earth, a retaining wall of the type noted above is plantable. Vegetation planted in those spaces decorates the retaining wall and may help to anchor the modules.\nRetaining walls of the type noted above and concrete modules therefor are exemplified in European Patent Application (Publication) No. 13,535, Steiner U.S. Pat. No. 4,379,659 and U.S. Pat. No. 4,521,138, Zeidman U.S. Pat. No. 5,072,566, Rausch U.S. Pat. No. 5,108,231, and Winkler et al. U.S. Pat. No. 5,177,925.\nAs exemplified in Swiss Patent No. 587,390 and corresponding German Published Patent Application (Offenlegungsschrift) No. 2,537,408, it is known for such a concrete module to have an open-topped, open-backed, trough-like configuration, which is defined by unitary front, lateral, and bottom walls. As illustrated and described therein, the front wall extends above the lateral walls and beyond the lateral walls on each side of the module, and the bottom wall is recessed upwardly from the lower edges of the front and lateral walls.\nAs a practical matter, the maximum height for a retaining wall built with such concrete modules, as described in the preceding paragraph, is limited by the columnar strength of the lateral walls of the modules. A maximum height of seven meters therefor is suggested by Winkler et al. U.S. Pat. No. 5,177,925, in column 1, lines 55 et seq. Such height limitations will be substantially more severe with curved (i.e., concave or convex) retaining walls.\nAs exemplified in Gavin U.S. Pat. No. 5,066,169, it is known to anchor a geogrid to concrete modules for a retaining wall, via retainer bars coacting with the geogrid and with integral bars on the modules.\nConcrete modules of related interest, some having means for anchoring geogrids thereto, are exemplified in Broadbent U.S. Pat. No. 4,470,728, Giardini U.S. Pat. No. 4,671,706, Forsberg U.S. Pat. No. 4,914,876, Miner U.S. Pat. No. 4,936,714, Rossi U.S. Pat. No. 4,964,761, and Janopaul, Jr., U.S. Pat. No. 5,044,834.\nRetaining walls of related interest, some employing geogrids, are exemplified in Broms et al. U.S. Pat. No. 3,925,994, Hilfilker U.S. Pat. No. 4,616,959, O'Neill U.S. Pat. No. 4,684,294, Brown U.S. Pat. No. 4,824,293, Wilson et al. U.S. Pat. No. 5,028,172, Risi et al. U.S. Pat. No. 5,064,313, and Strassil U.S. Pat. No. 5,108,232."}
{"text": "Field of the Invention\nEmbodiments of the present invention relate to variable nozzle turbochargers.\nDescription of the Related Art\nJapanese Utility Model Publication No. 63-115532 and Japanese Utility Model Publication No. 63-92036 disclose a variable nozzle turbocharger equipped with a variable nozzle mechanism. The variable nozzle mechanism adjusts the degree of opening of variable nozzles having nozzle vanes through rotation of a unison ring. Through this adjustment, the flow velocity of exhaust gas to a turbine wheel is controlled. In Japanese Utility Model Publication No. 63-115532, plate springs respectively biasing the variable nozzles in an axial direction of the unison ring are fixed to the unison ring. As a result, it is possible to prevent rattling caused by a gap between a shaft portion of each variable nozzle and a housing supporting the shaft portion. In Japanese Utility Model Publication No. 63-92036, a tension coil spring biasing each variable nozzle in the opening direction is provided between a link member of each variable nozzle and the unison ring. As a result, it is possible to prevent rattling caused by a gap between each link member and the unison ring.\nIn a general variable nozzle mechanism, exhaust gas passing between the nozzle vanes of adjacent variable nozzles imparts a force so as to open the nozzle vanes. The force is transmitted to the unison ring and an actuator for driving the unison ring. Each member of the power transmission route within the actuator as well as each variable nozzle receives a force with which it is pressed in one direction. As a result, it is possible to suppress rattling caused by a gap or backlash or the like between the members connected together in the power transmission route.\nAs the nozzle vanes open, the amount of exhaust gas acting on the nozzle vanes decreases. In some cases, it is impossible to obtain the force required to press each member in one direction, e.g., the rotational force acting on the nozzle vanes so as to open the nozzle vanes with exhaust gas. Conventionally, to prevent this, the pivot ratio (or the rotary shaft position) of the variable nozzle, or the configuration of the nozzle vanes, has been changed. This, however, results in an increase in the pressure loss of the exhaust gas, thereby resulting in performance deterioration.\nEach of the springs as disclosed in Japanese Utility Model Publication No. 63-115532 and Japanese Utility Model Publication No. 63-92036 is provided between each variable nozzle and the unison ring. The springs do not bias the unison ring so as to open the variable nozzles, but rather the springs bias the variable nozzle. Thus, the springs cannot cope with the above problem.\nTherefore, there is a need in the art for a variable nozzle turbocharger capable of preventing rattling caused by a gap, backlash or the like between members that are coupled together in the power transmission route of a variable nozzle mechanism."}
{"text": "Computer users may access the vast resources of the World Wide Web using web browsers and other like client-side applications, which generally operate by receiving code written in a computer language such as Hypertext Markup Language (or “HTML”) from one or more servers and/or server-side applications according to a certain protocol, such as Hypertext Transfer Protocol (or “HTTP”); parsing the code into a Document Object Model (or “DOM”); and assembling the DOM into a web page that may be displayed to a user on a user interface. Such browsers or other applications typically include engines for translating information received from a server into content that may be displayed on the user interface, interpreters for parsing and executing script, and databases or other data stores for storing information in the form of “cookies,” or files that relate to a user's browsing activity and may be retrieved by the same web site or other web sites in the future.\nThe display of a web page on a web browser begins with the initiation of a load event, which causes the web browser to transmit a request for one or more files associated with a Uniform Resource Locator (or “URL”) or other file location targeted by the load event. Such load events may include the entry of a URL into an address box or address bar; the selection of an image or set of text that may be hyperlinked to a URL; the selection of a bookmark, a home button or any other feature that may be linked to or otherwise associated with a URL; or a shut-down of the web browser and a subsequent return to the web page upon a restart of the web browser. When a user of a web browser provides the browser with a load event, the web browser may perform a search of various data stores, sometimes called a Domain Name System (or “DNS”) “lookup,” for an Internet Protocol (or “IP”) address associated with the URL. For example, when attempting to locate IP addresses that may be associated with the URL, the browser may search its own cached records, or records stored in caches on the operating system or router through which access to the Internet is provided, as well as external sources.\nOnce an IP address associated with the URL is identified, the browser sends a request for the web page at the URL, in accordance with the HTTP, to the server where the files associated with the page are located, and the server handles the request by providing a response to the browser in the form of code expressed in one or more computer languages, such as HTML. The browser may then begin to render the code into the web page, and may transmit requests, as necessary, for files that may be embedded within the code provided by the server. For example, where the code comprises HTML and refers to images, multimedia or other files that may be found at other locations (e.g., on external servers), the browser may transmit requests for each of the files. As the code and other files are received by the browser from the server, the browser may begin to assemble the web page for display on a user interface.\nOccasionally, during the loading of a web page, or after the loading of a web page has completed, a user of a web browser may elect to initiate a reload event which results in the retransmission of a request for the one or more files associated with the same URL that was provided in the load event that prompted the initial loading of the page. Most commonly, a reload event includes the selection of a “reload” or “refresh” button, which resubmits the URL for the web page that was most recently submitted to the servers, and causes the web browser to begin to load the web page again. Other reload events may include, but are not limited to, the reentry of the same URL that was provided with the load event into an address bar or address box, the reselection of an image or set of text that is hyperlinked to the same URL that was provided with the load event (i.e., selecting a hyperlink in a web page or electronic mail document, the selection of a bookmark for a web page that is currently displayed), or any step that returns the web browser to the web page following a shut-down or other action that secures the operation of the web browser.\nA user of a client-based application such as a web browser may initiate a reload event for a variety of reasons. First, the user may be dissatisfied with the rate at which the web page, or one or more elements expressed thereon, loads on the display. Second, one or more of the elements intended to be expressed on a web page, or the entire web page itself, may ultimately fail to load. Third, and perhaps most importantly, the user may be hindered or prevented from performing one or more tasks that prompted the user to visit the web page in the first instance. In this regard, the user's initiation of a reload event, which results in the submission of a second or subsequent request for a web page, may be detected and interpreted as an expression of dissatisfaction with at least one aspect of the web page, or an indication that the user is experiencing a difficulty in performing one or more tasks such as those that may ultimately be performed after the reload event. Accordingly, the initiation of a reload event, or the frequency with which users initiate reload events (e.g., the number of reload events initiated per unit time), may be tracked and analyzed to determine the interactivity of the web page and/or whether any remedial measures are required with respect to the web page, the client-side application or the server on which the web page and associated content is stored."}
{"text": "Technologies of recognizing and predicting feelings and behaviors of humans using machines are widely needed for society. For example, in order to realize robots more friendly with humans, it is necessary to recognize the feeling of other person. When whether or not a person behaves suspiciously can be determined from behavior prediction based on camera images from a monitoring camera, the number of staffs performing monitoring tasks can be reduced to save effort. An example of a key important to recognize the feelings of humans and to predict the behaviors thereof includes the line of sight of a person. This is because, though humans think or determine behaviors based on information obtained from external fields, most of the information is visual information.\nFor example, it is assumed that a camera is pointed at a driver of a vehicle and the line of sight of the driver is estimated from a camera image obtained using the pointed camera. In such a situation, it is assumed that it is found from the estimated line of sight of the driver that the driver intends to turn the vehicle to the left without seeing a motorcycle approaching from the left side of the vehicle. In this case, a behavior such as involved driving can be predicted. In addition, it is assumed that a camera is pointed at a viewer of a commercial display and the line of sight of the person standing in front of the display is estimated from a camera image obtained using the pointed camera. In such a situation, it is assumed that it is found from the estimated line of sight of the person that the point of gaze thereof is focused. In this case, it is possible to obtain information that the viewer shows an interest in information displayed at a position of the point of gaze.\nThe line of sight is a straight line in a three-dimensional space. Therefore, in order to define the line of sight, it is necessary to set two points in the three-dimensional space. Various methods involving defining the line of sight are expected. A first method is based on the idea that an eyeball is assumed as a structure in which a sphere forming a crystalline lens and a partial sphere forming a cornea portion are combined and a center of the cornea partial sphere and a center of a pupil are assumed as two points for defining the line of sight. The first method is normally called a “corneal reflection method” and disclosed in, for example, JP 07-055941 A (hereinafter referred to as “Patent Document 1”). According to the corneal reflection method, a rotation angle of the eyeball in a coordinate system fixed to a head is calculated based on a reflection light image (cornea reflection image or Purkinje image) observed in an eye region when the eyeball is irradiated with near-infrared illumination light, the center of the pupil, and an illumination relative position. A representation of a head posture in a world coordinate system which is estimated using any method and an eyeball rotation are combined to estimate the line of sight in the world coordinate system. Note that the “world coordinate system” means a coordinate system arbitrarily determined without depending on both a person and a camera.\nHereinafter, a three-dimensional coordinate system determined in association with a camera photographing a person is referred to as a “camera coordinate system”, and a coordinate system determined in association with the head of the person is referred to as a “head coordinate system”. The eyeball can reflect light from various light sources, and hence it is normally difficult to obtain a correspondence between an illumination light source and the cornea reflection image. Therefore, the corneal reflection method is effective in a case where the camera photographing the eye region and the illumination light source are significantly close to the person (for example, camera and illumination light source are incorporated in head mount display). However, the corneal reflection method has a problem that the person is required to wear the camera or the like and thus has inconvenience.\nOn the other hand, the second method is based on the idea that the eyeball is assumed as a single sphere and the center of the eyeball and the center of the pupil are assumed as two points for defining the line of sight. According to this idea (second method), the line of sight can be estimated without depending on positions of the camera and the illumination light source with respect to the person. Hereinafter, a combination of an eyeball central position and an eyeball radius is referred to as eyeball parameters.\nThe pupil is a feature point appearing on a surface of the eyeball, and hence a pupil central position can be detected from the camera image. However, the center of the eyeball is positioned in an inner portion of a human, and hence the eyeball central position cannot be directly detected from the camera image.\nVarious related technologies of estimating the eyeball central position have been known.\nAccording to a face image processing system described in, for example, JP 2004-504684 A (WO 02/009025 A1) (hereinafter referred to as “Patent Document 2”), during a calibration operation for estimating the eyeball central position, a person to be examined gazes on the point of gaze with known coordinates in a plurality of postures. The eyeball central position is estimated based on the fact that the line of sight passes through the set point of gaze in each of the postures.\nIn a system described in JP 2003-015816 A (hereinafter referred to as “Patent Document 3”), it is assumed that the center of the eyeball is positioned on a straight line which passes through a middle point between the inner corner of the eye and the outer corner of the eye and has a directional vector determined from a head posture."}
{"text": "Interleukin-1 (IL-1) plays an important role in the pathogenesis of several inflammatory disorders. Two proteins displaying IL-1 activity have been described, interleukin-1-alpha (IL-1.alpha.) and interleukin-1-beta (IL-1.beta.).\nAlthough the two proteins are products of distinct genes, they share 27-33% amino acid identity, and are known to interact with the same receptor and have similar biological activity. Each protein is proteolytically cleaved from its approximately 31 kDa precursor to its more active 17 kDa form. IL-1.alpha. precursor (preIL-1.alpha.) displays bioactivity, although less than that of its mature form. In contrast, IL-1.beta. precursor (preIL-1.beta.) displays no biological activity until processed to its mature form. Although both IL-1 molecules are secreted proteins, both lack signal peptides. The mechanism(s) of secretion have not been fully defined.\nOnly certain cell types process preIL-1.beta. and secrete IL-1.beta.. Monocytes and macrophages are the most efficient and prolific producers and secretors of IL-1.beta. known. Of the two IL-1 proteins synthesized and secreted following activation of monocytes and macrophages, IL-1.beta. is the more abundant.\nThe cellular processing of preIL-1.beta. to IL-1.beta. is mediated by the enzyme, IL-1.beta. Converting Enzyme (IL-1CE or ICE). ICE is expressed in vivo as a 45 kDa precursor molecule which is proteolytically processed in vivo into fragments of 20 and 10 kDa, that together comprise the active form of ICE.\nThe study of IL-1.beta. and the role it plays in the pathogenesis of several inflamatory disorders has been hampered by an insufficient amount of ICE. In vivo, monocytes and macrophages contain only small quantities of IL-1.beta. and ICE. Both ICE and IL-1.beta. are present in low quantities in cells, which limits the amount of material available for study and characterization. It has been reported that active ICE could not be produced by mixing the major protein subunits together, after separate expression in E. coli, or by coexpressing the subunits in Escherichia coli. Miller et al, \"Purification and Characterization of Active Human Interleukin-1.beta.-converting Enzyme from THP.1Monocyte Cells\", J. Biol. Chem. 268, pp 18062-18069 (1993). The present invention directly contradicts the teaching of Miller et al.\nSubstances that interact with ICE and alter the production of IL-1.beta. are of interest as therapeutics to modulate the inflammatory response."}
{"text": "The instant invention relates to toy vehicles and more particularly to a sparking toy vehicle and a launcher therefor.\nSparking toy vehicles and launchers therefor have heretofore been known in the art. In this regard, the U.S. Pat. Nos. to Kennedy No. 4,479,326 and Kakizaki No. 4,571,212 represent the closest prior art to the subject matter of the instant invention of which the applicant is aware. The patent to Kennedy discloses a toy vehicle projecting gun assembly in which the vehicle is normally latched on an inclined ramp with the rear wheels raised and free to turn. Joined to the underside of the gun is a power trigger operated by the trigger finger of the operator. On the upper side of the gun is a firing button operated by the user's thumb. Actuation of the trigger drives a gear train which in turn drives a drive gear mounted on the rear wheel axle, the drive gear being coupled to a fly-wheel. The flywheel is energized by repeated trigger actions, whereupon the player presses the firing button which unlatches the vehicle which then shoots down the ramp. The patent to Kakizaki discloses a sparking toy vehicle which is driven by a fly-wheel capable of being energized by pulling a rack gear across a pinion which is connected to the fly wheel. One of the side surfaces of the flywheel is provided with a material, which when contact by a flint, is operative for producing sparks. The flint is maintained in contact with the fly wheel surface by means of a spring-biased flint holder.\nThe instant invention provides a toy comprising a toy vehicle, a launching assembly, and a rotatable sparking mechanism for generating sparks before the vehicle is launched from the launching assembly. In a first embodiment, the launching assembly includes a base track having a rack gear, and a push handle including a spring biased plunger for engagement with the rear end of the vehicle. The sparking mechanism is mounted in the vehicle and a drive gear for sparking mechanism extends downwardly from the vehicle where it intermeshes with the rack gear when the vehicle is received on the base track. The push handle is slidably movable on the base track for advancing the vehicle along the base track wherein advancement of the vehicle causes rotation of the sparking mechanism thereby generating sparks, and further wherein a resistance of the sparking mechanism maintains the plunger in a depressed disposition until the drive gear disengages the rack gear. In a second embodiment, the rotatable sparking mechanism is built into the push handle wherein advancement of the push handle along the base track causes rotation of a drive gear on the push handle thereby generating sparks. The plunger is selectively maintained in a depressed position until the drive gear disengages from the rack gear. In a third embodiment, the launching assembly comprises an inclined launching surface, and a gear train including a gear extending upwardly through an aperture in the launching surface. The gear train is driven by a pivotable lever which is mounted at one end to the shaft of an actuator gear of the gear train. A cam is also mounted to the shaft of the actuator gear wherein pivoting of the lever causes the cam to engage, retract and release a spring plunger for propulsion of the car off the launcher.\nAccordingly, it is an object of the instant invention to provide a sparking toy vehicle.\nIt is another object to provide a launcher for a toy vehicle.\nIt is yet another object to provide a launcher for a sparking toy vehicle which utilizes resistance between the sparking mechanism and a gear train to maintain a launching spring in a compressed position.\nOther objects, features and advantages of the invention shall become apparent as the description thereof proceeds when considered in connection with the accompanying illustrative drawings."}
{"text": "Not applicable\nNot applicable\n1. Field of the Invention\nThe present invention relates to straps for flat bed trailers of 18 wheelers. More particularly, the present invention relates to a tool for assisting in winding of such straps.\n2. General Background of the Invention\nStraps for flat bed trailers of 18 wheelers are usually wound up by hand when they are not needed. This takes a fair amount of time. Also, the ratchets on which the straps are wound get hot in the summer and cold in the winter, and are thus often uncomfortable to handle. Further, sometimes the straps can freeze, which makes it really difficult to wind the ratchet by hand.\nThe following U.S. Patents are incorporated herein by reference:\nU.S. Pat. Nos. 4,390,141; 4,693,432; 4,913,608; 5,295,664; 5,426,826.\nThe apparatus of the present invention solves the problems confronted in the art in a simple and straightforward manner. What is provided is apparatus including a strap winding tool which includes a bent metal bar with a short elbow which can fit into a bore hole opening in a ratchet extension on a strap winder on a flat bed truck."}
{"text": "A wiper device having a nozzle unit, which is provided for arrangement on a wiper arm and for applying washing water to a vehicle window has already been proposed."}
{"text": "In the case of packet transmission, a packet transmission apparatus divides packet data, which are received from its preceding communication apparatus, into a plurality of blocks each having a fixed block size (predetermined block size), and assigns the blocks to a plurality of transmission lines. Accordingly, the packet transmission apparatus can increase transmission capacity by use of the plural transmission lines.\nIn the case that a size of the packet data, which are received from the communication apparatus, is not an integral multiple of the fixed data block size, if dividing the packet data into the blocks each having the fixed block size, n (n is any integer) blocks each having the fixed block size, and one block including residual data, whose size is smaller than the fixed block size, are generated.\nIn the case of a method according to a related art, by carrying out a padding process in which dummy data are added to the residual data, the packet transmission apparatus makes the size of the residual data adjusted to the fixed block size. Then, the packet transmission apparatus transfers data for which the padding process is carried out. The method has a problem that data transmission efficiency is lowered, since the transferred data includes a block which has an area unused for data transmission.\nIn order to solve the problem, PTL 1 discloses an art to set a size of packet data, which are sent by a communication apparatus, to be equal to or smaller than MTU (Maximum Transfer Unit) defined between the communication apparatus and a transmission destination apparatus, and to be an integral multiple of a block size which is applied to a low rank transmission line. By using the method, it is avoided to generate the residual data in the division process carried out by the packet transmission apparatus which receives the packet data from the communication apparatus. Therefore, it is possible to avoid generation of the block which includes the unused area."}
{"text": "The present application relates to semiconductor device fabrication. More particularly, the present application relates to a flash memory device that can be fabricated using a standard complimentary metal-oxide-semiconductor (CMOS) process.\nEmbedded non-volatile memory (NVM) devices have been widely used to store code and security information or trim information in integrated circuits. Among them, embedded flash memory cells in particular have been incorporated into CMOS integrated circuits. Commercially available flash memory cells generally include a control gate, a floating gate and two interposed dielectric layers disposed upon a semiconductor substrate. The presence of the additional floating gate in flash memory cells compared to field effect transistors requires dedicated device processing in order to integrate the flash memory cells into an integrated circuit. This includes additional steps of depositing an inter-gate dielectric and a control gate, raising additional challenges in gate patterning and gate spacer formation. These additional processing steps result in longer processing times, higher manufacturing costs and low yield. As such, a process flow to fabricate flash memory cells that is compatible with a standard CMOS fabrication process is highly desirable."}
{"text": "Demand for ever increasing realism and detail in computer-implemented video games seems to drive the growth of computer performance. Unlike computer animation and movie rendering which can process individual scenes over time for playback at a higher frame rate, computer-implemented video games and computer-implemented simulators must render complete, three-dimensional (3D) scenes of a virtual game environment during runtime of the game application, typically at a rate of thirty (30) frames per second or better. It can be difficult to produce animations that appear lifelike, fluid, and realistic when rendering during runtime of the game application."}
{"text": "U.S. Pat. No. 7,108,328 describes an armrest comprising an arm support rotatably mounted for pivoting about an axis and that can be locked in various positions by a latch movable between a locked position and a released position. A lockable ratchet mechanism is provided for this purpose, so that the arm support can be moved upward in the locked position, but a downward movement is blocked. The latch moves out of the locked position into the released position only when the arm support has been pivoted upward beyond a certain upper position, whereupon the arm support can be pivoted downwardly. As soon as the arm support is pivoted beyond a certain lower position, the latch moves back into the locked position.\nThe latch has two annular parts with interengaging toothed faces, namely a first part connected to the armrest in terms of movement and a second part nonrotatable about the axis. The second part is mounted on the axis so as to be displaceable axially and, in the locked position is engaged with the teeth of the first part and in the released position is out of mesh with the teeth of the first part. A spring force applies a load to the second part so as to mesh its teeth with those of the first part.\nThe pivot axis, the first part and the second part are formed of metal to be able to absorb high retaining forces despite having a low space requirement, so that the latch does not significantly impair the design of the armrest."}
{"text": "1. Field\nEmbodiments of the invention relate to electronics, and more particularly, to protection of communication transceivers from electrostatic discharge.\n2. Description of the Related Technology\nCertain electronic systems can be exposed to a transient electrical event, or an electrical signal of a relatively short duration having rapidly changing voltage and high power. Transient electrical events can include, for example, electrical overstress/electrostatic discharge (EOS/ESD) events arising from the abrupt release of charge from an object or person to an electronic system.\nTransient electrical events can damage integrated circuits (ICs) inside an electronic system due to overvoltage conditions and/or high levels of power dissipation over relatively small areas of the ICs. High power dissipation can increase IC temperature, and can lead to numerous problems, such as gate oxide punch-through, junction damage, metal damage, and surface charge accumulation. Moreover, transient electrical events can induce latch-up (in other words, inadvertent creation of a low-impedance path), thereby disrupting the functioning of the IC and potentially causing permanent damage to the IC."}
{"text": "Services can be provided in a computing environment. The environment can be a virtualized computing environment. The services can include infrastructure as a service (IaaS), platform as a service (PaaS), and software as a service (SaaS). End users can find and subscribe to these services via a portal. The services can be cataloged, monitored, and maintained to create a stable service offering."}
{"text": "1. Field of the Invention\nThe present invention relates to a head protecting airbag system that deploys an airbag downward along the side of a vehicle body in a passenger compartment when the vehicle is involved in either a side collision or a rollover. In particular, the present invention relates to a head protecting airbag system that deploys an inflation section between an occupant's head and the vehicle body.\n2. Description of the Related Art\nIn the recent trends, a head protecting airbag system is mounted in vehicles as a supplemental restraint system. The head protection airbag system has a curtain airbag that is deployed downward from a roof side rail when the vehicle is involved in a side collision or rollover. As one of such head protecting airbag systems, there has been proposed a head protecting airbag system designed to deploy the airbag between the head of the occupant and the side of the vehicle body in order to protect the occupant's head.\nThis type of head protecting airbag system is described in Japanese Patent Application Publication No. 2004-58848 (JP-A-2004-58848). In this head protecting airbag system, the airbag includes a gas supply path and a main inflation part disposed on the lower side of the gas supply path. The airbag is folded with the main inflation part rolled up toward the outer side in the vehicle and with the gas supply path not rolled, but folded for the sake of easier deployment of the airbag upon gas supply. In addition, in Japanese Patent Publication No. 3656156, at the initial stage of deployment of the airbag, a part of the airbag deploys into a thin plate shape between the pillar and the headrest of the occupant seat with no gas supplied. Then, gas is supplied into the thin plate airbag from an inflow port formed on the lower side of the airbag. This prevents the inflating airbag from being interfered with by the headrest of the occupant seat.\nHowever, in order to immediately deploy the airbag in the longitudinal direction of the vehicle along the side of the passenger compartment, the airbag has to initially deploy inward towards the passenger compartment over the top end of the pillar garnish."}
{"text": "1. Field of the Invention\nThe present invention relates to a medical image display control device, a method of operation for the same, and a medical image display control program for displaying the internal tissue of a subject so as to be superimposed on a three-dimensional image of the subject.\n2. Description of the Related Art\nIn the related art, the planning of actual surgery or the like has been performed using a three-dimensional image captured before surgery. When performing such surgical planning, the internal tissue, such as blood vessels at the back of the skin tissue or the bone tissue, needs to be visualized in some cases.\nSpecifically, for example, in the case of performing craniotomy, only things on the surface are visible when performing actual surgery since the head has a layer structure of a skin region, a bone region, and a brain region from the surface. Therefore, when incising the skin, blood vessels under the skin are the important tissue for the surgery. In addition, when opening the bone, brain veins or brain arteries under the bone are the important tissue for the surgery.\nIn such a case, in general, the internal tissue, such as blood vessels, can be displayed by extracting the skin tissue or the bone tissue and adjusting the opacity or the attenuation.\nIn order to display the internal tissue, JP2007-325629A has proposed to display a three-dimensional image of the internal tissue excluding, for example, a region of the skin tissue by setting a predetermined range from the surface of the subject as a removal layer and constructing a three-dimensional image from which the removal layer has been removed."}
{"text": "1. Field of the Invention\nThis invention relates generally to shaft seals that are used in hydrogen inner-cooled turbine-generators, and more particularly to methods and apparatus for cooling the shaft seals in such turbine-generators.\n2. Statement of the Prior Art\nConductor cooling is a conventional process used in very large turbine-generator systems for dissipating the armature and field coil losses of the turbine-generators to cooling media within their coil insulation wall. The turbine-generators using such conductor cooling are also referred to variously as \"inner-cooled,\" \"supercharged,\" or \"direct-cooled\", and the cooling medium that is most often used in such turbine-generators is hydrogen.\nHydrogen inner-cooled turbine-generators typically operate in a pressurized hydrogen atmosphere that provides cooling for all of the turbine-generator except, in some instances, its armature coils. Such hydrogen inner-cooled turbine-generators are often operated at 60 lbf/in.sup.2 (i.e., about 4219 grams per square centimeter) or more in order to increase the mass flow of the hydrogen, and to reduce its temperature rise.\nIn order to minimize leakage of pressurized hydrogen cooling medium from an operational turbine-generator, shaft seals are typically used in hydrogen inner-cooled turbine-generators for maintaining an oil film under pressure in a small clearance between the rotating shaft of the turbine-generator and a stationary member surrounding the shaft at both ends of the turbine-generator. The construction of such shaft seals may be similar to a journal bearing with a cylindrical oil film or similar to a spring-loaded thrust bearing with the oil film in a plane at right angles to the shaft axis. In either case, the oil film is maintained by an oil supply pressure that is higher than the hydrogen pressure.\nOils used in such shaft seals can absorb about 10% by volume of either hydrogen or air. It is important that the flow of oil in those shaft seals toward their hydrogen side be minimized in order to reduce both the amount of air that is carried into the hydrogen inner-cooled turbine-generator and the amount of hydrogen that is carried out. Moreover, it is important to minimize temperature differences between the \"hydrogen side\" of the shaft seal and its \"air side\" so that differential thermal expansion of the seal ring which comprises the stationary member can be minimized.\nPrior art shaft seals have typically employed separate supplies of oil for their hydrogen side and their air side, each such oil supply including a heat exchanger operated in one of two general fashions. One method selects a cooling media (e.g., water), and carefully sizes the heat exchanger of each oil supply to provide for open loop control of seal oil temperatures at the outlets on either side of the shaft seals. Any problems with heat exchanger fouling, decreased oil flow or decreased water flow when experienced with this method requires recognition and subsequent correction by an operator.\nAnother method utilizing separate heat exchangers for the hydrogen side and the air side of prior art shaft seals provides automated control valves to sense the temperature difference between the hydrogen side oil supply and the air side oil supply, and operates the particular heat exchanger in the oil supply having the higher temperature. Not only is this other method complicated in its design by virtue of the additional components that are required (e.g., two heat exchangers and many automated control valves), but it is also difficult in its implementation due to the necessity for precise, reliable components to accomplish such control over a very narrow temperature range."}
{"text": "1. Field of the Invention\nThe present invention relates to an information providing apparatus in which unspecified many people can arbitrarily use it to obtain various kinds of information with a fee or without a fee.\n2. Description of the Related Art\nIn recent years, various kinds of information have been provided in an electronic information mode in which they are recorded in a medium.\nFor example, a newspaper, a magazine, a novel, and the like are recorded in CD-ROMs or the like as electronic publications and are on sale. The same are applied to music, video (movie, video clip, etc.), game software, application software used in a personal computer, and the like, and various kinds of information (hereinafter, information recorded in a medium and provided will be generically referred to as xe2x80x9ccontentxe2x80x9d in the present specification) are provided with a fee or without a fee in an extremely wide range.\nIn the content which is electronized and is provided, there is content which a user keeps for a long period (music software, application software, etc.) and content which the user feels that long storage is unnecessary.\nFor example, when it is considered that a newspaper or a weekly magazine is electronized and provided, if the user sees the information once, it often becomes unnecessary thereafter.\nIn other words, it is not realistic to sell a newspaper or a magazine, as a package medium, in a keepable recording medium of a disk or card. For example, if the user sees the contents once, not only the information but also the medium itself becomes unnecessary, so that useless consumption of resources is caused.\nOn the other hand, with respect to the information which does not require long storage, it is also carried out for a user to search or download arbitrarily from his own terminal (personal computer, etc.) through communication such as the Internet. However, in this case, of course, in order to see the contents of the information, the user must be in a place (home, office, etc.) where the terminal is set. For example, in the case where the user desires to see the contents of the information at an arbitrary place through a portable disk reproducing apparatus or the like, it is necessary to further download the information from the terminal to a disk medium or the like.\nBesides, in order to obtain content such as an electrical newspaper through communication, it is naturally necessary for the user to possess the terminal connected to a communication network.\nFrom the above circumstances, it has been required to realize such that a general user can easily obtain a medium which records content such as an electronic newspaper or magazine, and the user can confirm the content at an arbitrary place by using, for example, a portable reproducing apparatus or a home terminal, and further, waste of the medium itself does not occur.\nBesides, when considering an apparatus in which a general user can obtain a medium which records content, it is also required that inventory control can be simply made.\nIn view of the foregoing circumstances, an object of the present invention is to provide a content providing system in which unspecified many people arbitrarily use it to obtain (purchase) a recording medium in which desired content has been downloaded, and the recording medium itself can be effectively used, and further, inventory control of the recording medium recording the content to be sold can be made simple and adequate.\nAccording to an aspect of the invention, an information providing apparatus comprises: a loading section in which a recording medium including use state identification information indicating one of a personal use and an unspecified user\"\"s use is loaded; recording means for recording content information in the recording medium loaded in the loading section; use state detecting means for detecting which of the personal use and the unspecified user\"\"s use the use state identification information of the recording medium loaded in the loading section indicates; ejecting means for ejecting the recording medium to an outside of the apparatus; content detecting means for detecting a kind of content information to be recorded in the recording medium ejected by the ejecting means; a housing section capable of housing a recording medium in which one of various kinds of content information is recorded in advance; and control means, wherein when the use state detecting means detects that the use state identification information of the recording medium indicates the personal use, the control means controls the recording means so that the content information of the kind detected by the content detecting means is recorded in the recording medium loaded in the loading section, and then, the control means controls the ejecting means so that the recording medium loaded in the loading section is ejected; and when the use state detecting means detects that the use state identification information of the recording medium indicates the unspecified user\"\"s use, the control means controls the ejecting means so that the recording medium in which the content information of the kind detected by the content detecting means is recorded is ejected among recording media housed in the housing section.\nAccording to another aspect of the invention, an information providing apparatus comprises: a housing section capable of housing a plurality of recording media in which any one of plural kinds of contents information is recorded in advance; content detecting means for detecting a kind of content information required by a user; ejecting means for ejecting a recording medium in which the content information of the kind detected by the content detecting means among recording media housed in the housing section; remaining recording medium detecting means for detecting whether or not the remaining number of the recording media which are housed in the housing section and in which each kind of content information is recorded, is a minimum threshold value or less; and replenishing means for replenishing the recording medium in which the content information of the kind detected by the remaining recording medium detecting means such that the remaining number is the minimum threshold value or less.\nAccording to further aspect of the invention, an information providing apparatus comprises: a loading section in which a recording medium is loaded; content identification information detecting means for detecting content identification information of the recording medium loaded in the loading section; setting means for setting content information on the basis of the content identification information detected by the content identification information detecting means; housing section capable of housing a plurality of recording media in which content information indicated by the content identification information is recorded in advance; and ejecting means for ejecting, among recording media housed in the housing section, a recording medium in which the content information coincident with the content information set by the setting means is recorded.\nAccording to still further aspect of the invention, a recording medium which loaded in an information providing apparatus in a removable manner and in which content information can be recorded by the information providing apparatus, wherein the recording medium comprises: a recording medium main body; and a content identification section including content identification information indicating a kind of content information which can be detected by the information providing apparatus.\nAccording to still another aspect of the invention, a method of providing content comprises the steps of: detecting which of a personal use and an unspecified user\"\"s use is indicated by use state identification information of a recording medium loaded in a loading section; detecting a kind of content information to be recorded in a recording medium ejected by ejecting means; wherein when it is detected in the step of detecting a use state that the use state identification information of the recording medium indicates the personal use, recording the content information of the kind detected in the step of detecting the kind of the content information in the recording medium loaded in the loading section, and then, controlling the ejecting means so that the recording medium loaded in the loading section is ejected; and when it is detected in the step of detecting the use state that the use state identification information of the recording medium indicates the unspecified user\"\"s use, controlling the ejecting means so that, among recording media housed in a housing section capable of housing recording media in which any one of various kinds of content information is recorded in advance, the recording medium in which the content information of the same kind as the kind detected in the step of detecting the kind of the content information is recorded, is ejected.\nAccording to still another aspect of the invention, a method of providing content comprises: step (a) of detecting a kind of content information required by a user; step (b) of ejecting a recording medium in which the content information of the kind detected in the step (b) is recorded, among recording media housed in a housing section capable of housing a plurality of recording media in which any one of various kinds of content information is recorded in advance; step (c) of detecting whether or not the remaining number of recording media which are housed in the housing section and in which each kind of content information is recorded is a minimum threshold value or less; and step (d) of replenishing the recording medium in which the content information of the kind which has been detected in the step (c) that the remaining number is the minimum threshold value or less, is recorded.\nAccording to still another aspect of the invention, a method of providing content comprises the steps of: detecting content identification information of a recording medium loaded in a loading section; setting content information on the basis of the content identification information detected by the step of detecting the content identification information; and ejecting a recording medium in which content information coincident with the content information set at the setting step is recorded, among recording media housed in a housing section capable of housing a plurality of recording media in which the content information indicated by the content identification information is recorded in advance."}
{"text": "1. Field of the Invention\nThe present invention relates to a fuel conditioner adapted to be placed in-line in a fuel delivery system and is adapted to condition fuel for improved purity, extended storage life, and reduced-engine wear.\n2. Description of the Related Art\nInternal combustion engine systems typically are provided fuel from a remote storage tank via fuel lines and the fuel is driven either by gravity or an active pump. The systems often include an in-line filter to remove particulate impurities. However, filters typically are passive devices that can only screen out particles above a certain size.\nFilters are relatively ineffectual against biological processes that often occur in fuel. As an example, diesel fuel often accumulates water in storage. The fuel, especially with water present, can support the growth of certain bacteria, fungi, and algae. It is known to add anti-biological agents such as bactericides and fungicides to the fuel, however this requires the undesirable additional effort of adding the anti-biological agents to the fuel.\nAn additional contaminant that can be present in fuel that is not particularly well handled by conventional filters is metal contamination. Small particles of metal can become entrained in the fuel from wear in fuel pumps and corrosion in fuel delivery systems. These small metallic particles can be too small to be effectively trapped by a filter, yet large enough to cause undesirable wear and deposits in the engine.\nFrom the foregoing, it can be appreciated that there is a need for a fuel conditioning system that can inhibit the growth of biological contaminants in fuel yet avoids the inconvenience of mixing additives with the fuel. There is also a need for a system to remove metallic particles from a fuel supply.\nThe aforementioned needs are satisfied by the present invention which, in one aspect, is an in-line fuel conditioner receiving a flow of liquid fuel the conditioner comprising a housing, end caps attached to ends of the housing and a plurality of reactive elements contained within the housing such that the fuel passes over the reactive elements and wherein the reactive elements comprise separate stainless steel, zinc, and copper members such that the overall composition of the reactive elements is approximately 50-40% stainless steel, 40-30% zinc, and 30-20% copper by weight and wherein the in-line fuel conditioner inhibits the growth of biological agents entrained within the fuel. In certain embodiments, the invention also includes a magnet wherein the magnet retains ferrous metal particles entrained within the fuel flow and in one embodiment, the magnet is positioned with the housing. In certain embodiments, the reactive elements are approximately 0.125xe2x80x3 in major dimension. In another aspect, the invention is an internal combustion engine system utilizing fuel and including an in-line fuel conditioner wherein the in-line fuel conditioner contains a plurality of reactive elements comprising separate stainless steel, zinc, and copper members wherein the reactive elements have a weight composition of 50-40% stainless steel, 40-30% zinc, and 30-20% copper and a magnet wherein the in-line fuel conditioner inhibits the growth of biological contaminants and retains ferrous particulates entrained within the fuel.\nThe in-line fuel conditioner of the present invention can be readily installed in an existing fuel delivery system using commonly available tools and known mechanical techniques. The fuel conditioner inhibits the growth of biological contaminants without the inconvenience of treating the fuel with additives. The fuel conditioner also retains ferrous particulates entrained within the fuel thereby reducing wear to an engine system so equipped. These and other objects and advantages will become more fully apparent from the following description taken in conjunction with the accompanying drawings."}
{"text": "Digital video services such as transmitting digital video information over wireless transmission networks, digital satellite services, streaming video over the internet, delivering video content to personal digital assistants or cellular phones, etc., are increasing in popularity. Increasingly, digital video compression and decompression techniques may be implemented that balance visual fidelity with compression levels to allow efficient transmission and storage of digital video content."}
{"text": "Java Platform, Enterprise Edition (Java EE) (the latest version of which is known as Java EE 5—the previous version was known as Java 2 Enterprise Edition, abbreviated to J2EE) is a set of coordinated technologies and practices that enable solutions for developing, deploying, and managing multi-tier, server-centric applications (especially distributed applications in which different components of the overall application run separately from one another). Building on Java Platform, Standard Edition (Java SE), Java EE adds the capabilities that provide a complete, stable, secure, and fast Java platform for the enterprise. Java EE significantly reduces the cost and complexity of developing and deploying multi-tier solutions, resulting in services that can be rapidly deployed and easily enhanced.\nNote that throughout this specification the term object will be used to refer to a computer program construct having various properties well known to persons skilled in the art of object oriented programming of which Java is a well known and commercially important example. Although the described embodiments of the present invention relate specifically to Java, it will be apparent that the invention is more generally applicable, especially as regards other object oriented programming languages.\nIn essence, Java EE enables an enterprise to develop a server application without needing to generate complex software for handling many of the functions generally required by server applications (e.g. of enabling connections to be made to the server application from remote clients, of providing security features to restrict access to features which only an administrator of the system should have access, etc.).\nIn particular, an enterprise may develop a server application by generating a number of small software modules known as “Enterprise Java Beans (EJBs)” which are fairly specific to the particular application. The EJB's are then “placed” into an EJB container which is an environment for running EJBs. The EJB container takes care of most of the low-level requirements of the application (e.g. implementing security and enabling remote clients to access the application).\nA number of different types of EJB are typically used to generate an application. There are three main different types of enterprise bean: Session beans, Entity beans and Message driven beans.\nSession beans are generally used to implement “business logic”—i.e. functions such as looking up data from a database and manipulating it to generate an output which is provided to a remote client. Session beans can either be stateless or “stateful”. Stateless session beans are distributed objects that do not have state associated with them thus allowing concurrent access to the bean. The contents of instance variables are not guaranteed to be preserved across method calls. Stateful session beans are distributed objects having state. The state can be persisted, and access to a single instance of the bean is limited to only one client.\nEntity beans are distributed objects having persistent state. They are generally used to represent an item of business specific data, especially data about a specific business entity (e.g. a piece of network equipment). In general, the persistent state may or may not be managed by the bean itself. Beans in which their container manages the persistent state are generally said to be using Container-Managed Persistence (CMP), whereas beans that manage their own state are said to be using Bean-Managed Persistence (BMP). In both cases the “persistence” is generally performed using a relational database as the backend data store and using an Object to Relational Mapping (ORM) function to convert between the object(s) associated with the entity beans and the backend relational database store. A particularly oft used ORM is Hibernate (see www.hibernate.org for more details of this product).\nMessage Driven Beans are distributed objects that respond to Java Messaging Service (JMS) messages. JMS is a standard part of the Java EE platform which enables java objects (typically on remote devices) to asynchronously communicate with one another by sending messages between them. Message beans were added in the EJB 2.0 specification to allow event-driven beans.\nA distributed server application may also include a type of object called a Managed bean, or Mbean. These are separate things to EJBs and are typically used to represent an application or device which an administrator would like to be able to remotely manage while the device and/or application is running.\nSince the rapid uptake of the internet and the overlying World Wide Web (WWW), the paradigm of distributed processing has been taken from a favourable method, to one that now tackles large management problems in large networks and a common method of designing applications. The ability to distribute a system allows for a delegation of resources, processing and even operational logic (that can be regarded as black box functionality) to other components and also allows duplication of components in different locations to provide redundancy.\nWith the uptake of the distributed paradigm though comes an increasing dependency; the lack of availability of a component can cost businesses customers, time and money. This means that components have to be readily available, redundant to allow for load and reliability, and perhaps most importantly; externally constant. The necessity for constancy stems from the lack of knowledge regarding systems that are dependent upon the component. Across a business-to-business (B2B) interface only the integration expert that originally coupled the two systems knows the effects of changing either of the two communicating components, and in his absence the upgrade of either component will be shied away from (opting for legacy support instead), or will require extensive and expensive analysis and then almost certainly the requirement for another adapter, B2B interface or mid-way component to be constructed. Such additions solve the immediate problem, but tend to obfuscate further the interaction between the two components should a further change be required.\nBut in the modern environment, particularly with web-based systems, there is so much change and redefinition that was not picked up with initial requirement analysis, in such a short period of time, that updates are not only necessary, but required increasingly more frequently.\nThis is true not only for wide scale (Wide Area Network (WAN) based) distributed applications, but also on same-host locally distributed applications. The popularity of container-based logic, brought to popular use by JAVA EE-type servers, uses a distributed-type paradigm in a local sense in order to separate systems' business logic into a series of process components that have instances managed by a container, and interrelate through message passing. This type of distribution is just as vulnerable to the dangers of updating; a single component unavailable can render the system unavailable or worse still erroneous. Although there should be more knowledge of the systems (i.e. all the code should still be available for the system), external interactions can still introduce data to the system that can stress it in unpredicted ways. Adding or updating components themselves still leaves downtime (during the replacement of a component) and could cause loss of instances or a processing chain if there were active instances running prior to the update. Varying from server type to server type, it can be very difficult to just deploy one component of an application at a time; there is no uniform standard in JAVA EE for component replacement or indeed how to deploy an application.\nCurrent methodologies to try and maintain “zero downtime” when updating or modifying a distributed application include server clustering of JAVA EE resources and federated applications. Both revolve around the concept of there being several servers running the same service or application. When an update is needed to be made, a server by server migration is conducted until all the hosts for the application are running the new version. This does allow for deployment analysis, and ensures that at no point is the service unavailable to the clients. What it doesn't ensure is ability to revert the logic after the update is complete, should an error be found. Also, although federating across multiple servers guarantees “zero down-time”, the propagation time for the new logic could be sufficient to allow for a second update to occur during the period, which could lead to significant versioning issues. Furthermore, it generally requires the enterprise to have a number of physical computers, which can be expensive for small enterprises.\nClustering and federation do not take any further steps to allow for integration or extension of data types within the system. Receiving classes that have been previously unspecified is possible, and they can be handled, however should these then need to be persisted they would have to undergo a transformation to a known similar class (which is no simple task and is likely to involve a loss of data) or cannot be stored at all (due to the possible lack of presence in the received class in the system libraries)."}
{"text": "1. Field\nEmbodiments of the invention relate to a multi-purpose tool which may include a saw, axe, mallet, and stake puller, and to a tension mechanism therefor.\n2. Description of the Related Art\nOutdoor tools can be used for a variety of tasks, such as cutting, sawing, chopping, woodworking, building, etc. One common outdoor tool is a bow saw used for sawing a workpiece such as a tree branch. Bow saws generally have a metal frame in the shape of a bow and a coarse wide blade attached to a handle. Another common tool is an axe, which is generally used to shape, split and cut wood or to harvest timber, for example. An axe may have many forms and uses, but generally includes an axe head with a handle, or helve. Another common tool is a mallet, which may be considered as a kind of hammer. Mallet heads are relatively large and are usually made of rubber, but can sometimes be made of wood, for example. A stake puller, as suggested by its name, is an outdoor tool that can be used to pull wooden or metal stakes from the ground."}
{"text": "1. Field of the Invention\nThe present invention relates to a projector for projecting an image onto a screen, and more particularly, to a projector having an improved structure for cooling an optical system thereof.\n2. Description of the Related Art\nA projector is a device generally including a light source and an optical system projecting and enlarging an image onto a screen, and is divided into two types, a CRT (Cathode Ray Tube) type and an LCD (Liquid Crystal Display) type.\nThe CRT-type projector (hereinafter referred to as a “CRT projector”) comprises a CRT assembly projecting an image beam, a mirror reflecting the projected beam, and a screen on which the reflected beam is displayed as a picture.\nThe LCD-type projector (hereinafter referred to as an “LCD projector”) comprises an LCD instead of the CRT assembly, wherein the image beam is adjusted in gradation by a liquid crystal and then projected and enlarged.\nSuch an LCD-type projector is divided into a single panel type and a triple panel type according to the number of LCD panels, and is advantageously brighter and lighter than the CRT projector.\nHereinafter, the triple-type LCD projector will be exemplarily described.\nGenerally, a high-power lamp is used to enhance brightness of a picture projected from the projector, and therefore a relatively large amount of heat is generated inside the projector. This heat may cause components inside the optical system to be heated and damaged. Accordingly, to manage temperature inside the projector, a forced convection cooling method using a cooling fan is employed.\nAs shown in FIGS. 1 and 2, a structure cooling an optical system 120 of a conventional projector comprises a duct 160 including an inlet 166 through which cool air flows in. The cooling structure also includes a plurality of outlets 167 opening toward an optical synthesizing part 130 including panels 131, 132 and 133 and target optical components 140 respectively. The cooling structure further comprises a cooling fan 150 placed inside the duct 160 that forcefully blows the cool air convectively.\nThe optical system 120 of the conventional projector operates as follows.\nThe cooling fan 150 forcefully blows the cool air, flowing in through the inlet 166 of the duct 160, to the outlets 167 and toward places where the heat is generated in the optical system 120, thereby cooling the optical system 120.\nFor example, some of the places where the heat is generated in the optical system 120, include the optical synthesizing part 130 including a blue panel 131, a green panel 132 and a red panel 133, and the optical components 140 including a flyeye-lens 141.\nHowever, color signals respectively transmitted to the blue panel 131, the green panel 132 and the red panel 133 vary in wavelength, so that respective frequencies and energies thereof are different from each other. For example, the blue panel 131 which processes a blue component of the image beam receives a color signals having the shortest wavelength and the highest frequency among the color signals, and therefore, the blue panel 131 generates maximum heat as compared with other panels.\nGenerally, the optical system 120 is cooled with reference to a temperature of the blue panel 131 which has the highest temperature among other panels, so that the red panel 133 and the flyeye-lens 141 which have a relatively low temperature are overcooled, thereby wasting available energy.\nFurther, when the optical system 120 is cooled with reference to the place having a relatively high temperature, the cooling fan 150 is driven fast so as to speed up air flow, thereby causing noise due to the rotation of the cooling fan 150 and the air flow."}
{"text": "In an active matrix liquid crystal display (LCD), an image is produced by controlling the light transmittance of a two-dimensional array of discrete image elements (subpixels). Control is performed by the conversion of digital image data, consisting of a data value for each subpixel of the image, into analogue voltages with values dependent on that data, and direction of those voltages to each pixel electrode in the array via an active matrix of source data lines, gate lines and thin film transistor (TFT) switching elements.\nA block of three subpixels is termed a pixel. Each subpixel is associated with a color filter (typically Red, Green and Blue), and by controlling the amount of light being transmitted through these three color channels, any resultant linear combination of Red, Green and Blue light can be produced. In the case of one or more color channels being in the “off” state—that is, the liquid crystals are energized in such a way that no light is transmitted through them—the color is said to be “saturated” in that the color is as vivid as is possible with this display. Some LCDs are capable of producing very vivid saturated colors, while some are only able to produce pale colors, even with only one color channel in the on state. The vividness of the color is related to its light spectrum. A light with a very broad spectrum will appear pale, while a narrow spectral light will appear vivid. A monochromatic light (that is, light with only one wavelength, typically produced by a laser) is the most vivid light possible.\nIn conventional LCDs, a broad spectrum light is emitted by a backlight, and each color filter will only transmit a particular range of wavelengths. A narrow range will produce more vivid colors, but since more light is being absorbed by the filter, the brightness of the screen is reduced. In addition, LEDs with a broad spectral emission are in general more power-efficient than those with narrower emission windows.\nThere have been several methods tried to improve the power efficiency of high-color screens. In Akiyama, U.S. Pat. No. 7,106,276 (issued Sep. 12, 2006) the inventors supplement the white LED (W LED) with three separate LEDs each with a narrow spectral emission, in the three color primaries. However, the power requirement of this configuration is large, and outweighs the advantage of the high-color display. Bergquist, Publication US20080150864 (published Jun. 26, 2008) specifies using only Red Green and Blue LEDs (RGB LEDs), in such a way that for each color primary, if the image content allows it, the corresponding LED can be dimmed and the image data altered. In this way, the power requirement can be reduced while maintaining the color vividness. However, this still requires the RGB LEDs to be used even for pale colors, which is less efficient than using W LEDs.\nVan Beek et al., US20090160756 (published Jun. 25, 2009) attempts to deal with this by combining RGB LEDs and W LEDs, and selectively choosing which LEDs to use at any one time. The inventors calculate the required drive current to each of the independently controllable R,G,B LEDs, and then make assumptions about the possibility of replacing R,G,B specific currents with a general W current. In this way they reduce the power consumption of the backlight device, but they also restrict the vividness of the image data to be displayed. Specifically, if their algorithm concludes that the Red, Green and Blue LEDs should all be driven at maximum current, then they will replace this with a White LED being driven at full current and the RGB LED not driven at all. This clearly reduces the vividness of the color that the panel can produce, even though the image content might require high saturation levels. More generally, for any color which fits inside the gamut of the White LED, this approach will dim the W LED to the smallest of the three colors, and supplement it with RGB LEDs, instead of the more efficient approach of using the W LED to supply all the required light, and not using the RGB LEDs.\nLangendijk et al., U.S. Pat. No. 8,300,069 (issued on Oct. 30, 2012) discloses another means of balancing different backlights. The inventors do not use a W LED, but use a fourth (White) subpixel. They control the current going to the Red LED, and to the Green and Blue LED. To avoid “chromaticity dependence”—that is, Green light being transmitted through the Blue color filter and vice versa—the ratio of G:B current is held constant. For pale pixels, the White subpixel can be opened in such a way that the pixel will appear very bright; for vivid pixels, the White will be closed and only the desired primaries will be transmissive. In this way, the apparent brightness of the panel can be increased in pale areas so that the overall brightness appears greater; or the backlight powers can be scaled down to reduce the power increase necessitated by having the RGB LEDs. However, this gives a non-uniform distribution of brightnesses, contrary to standardized color spaces.\nMany conventional devices (e.g., Nakano et al., U.S. Pat. No. 7,333,165 (issued Feb. 19, 2008), Keh et al., US20070103934 (published May 10, 2007), and Morishita, US20120242564 (published Sep. 27, 2012)) use a combination of different backlights, but the proportion of each backlight being used depends on a variety of different system conditions. For example, the ambient lighting, the desired brightness, or the display mode being used (e.g., Vivid, Low power), can all affect the balance."}
{"text": "Many health care procedures involve medication administrations. The type of medication and timing of administration are important to record in order to provide healthcare providers real-time information on the conduct of the procedure and the completion of a medical record. Some protocols require quick medication administrations with limited time for documentation and record keeping. Others require completion and verification of medication administration manually to ensure proper patient care and accounting for use of medications. To ensure proper safety, all medication administrations require the administering clinician to verify the “5 rights”—right patient, right medication, right dose, right time, and right route, which necessitates the collection, association, and recording of multiple data elements for each drug given. Additionally, medication administration requires careful management of injection sites to maintain patency and infection control."}
{"text": "1. Field of the Invention\nThe present application relates generally to an improved data processing system and method. More specifically, the present application is directed to a system and method for zoning of devices, such as serial attached SCSI (SAS) devices, for example, in a storage area network (SAN) based on logical unit number (LUN) masking/mapping.\n2. Background of the Invention\nStorage area networks, or SANs, consist of multiple storage devices connected by one or more fabrics. Storage devices can be of two types: host systems that access data, and storage subsystems that are providers of data. Zoning is a network-layer access control mechanism that dictates which storage subsystems are visible to which host systems. This access control mechanism is useful in scenarios where the storage area network is shared across multiple administrative or functional domains. Such scenarios are common in large installations of storage area networks, such as those found in storage service providers.\nThe current approach to zoning storage area networks is typically manual and involves correlating information from multiple sources to achieve the desired results. For example, if a system administrator wants to put multiple storage devices in one zone, the system administrator has to identify all the ports belonging to the storage devices, verify the fabric connectivity of these storage devices to determine the intermediate switch ports and input all this assembled information into the zone configuration utility provided by the fabric manufacturer. This manual process is very error-prone because storage device or switch ports are identified by a 48-byte hexadecimal notation that is not easy to remember or manipulate. Furthermore, the system administrator has to also do a manual translation of any zoning policy to determine the number of zones as well as the assignment of storage devices to zones.\nRecently a new approach to zoning of Fibre Channel architecture based storage area networks has been devised in which automatic zoning of storage area networks (SANs) is performed based on system administrator defined policies. This mechanism is described in commonly assigned U.S. Patent Application Publication No. 2005/0091353 entitled “System and Method for Autonomically Zoning Storage Area Networks Based on Policy Requirements.” With this mechanism, in a measurement phase of operation, configuration data and connectivity topology data is periodically collected from each of the devices in the SAN using software agents which are placed in every switch and host device in the SAN. The collected data is analyzed to determine relationships between the devices in the SAN. The information obtained from such an analysis is converted into a graph structure where each node is either a switch port or a storage device port. A zone plan is then generated based on the graph structure and system administrator established zoning policies. The generated zone plan is then provided to a system administrator for approval and, if approved, the zone plan is implemented in the SAN by programming the zoning plan onto individual switches included within the SAN.\nOne emerging technology that is being used more prevalently with storage area networks is Serial Attached SCSI (SAS) communication protocol technology. SAS is a computer bus technology primarily designed for transfer of data to and from devices such as hard drives, CD-ROM drives, tape storage devices, and the like. Traditionally, SAS is a serial communication protocol for direct attached storage (DAS) devices. Recently, however, the protocol has been extended to include storage area network support. It is designed for the corporate and enterprise market as a replacement for parallel SCSI, allowing for much higher speed data transfers than previously available, and is backwards-compatible with SATA drives. Though SAS uses serial communication instead of the parallel method found in traditional SCSI devices, it still uses SCSI commands for interacting with SAS end devices.\nSAS zoning in a storage area network is particularly difficult. Typically, the SAS zoning is a two step mapping where each individual port, or “phy,” is assigned to a zone group and then, as a subsequent step, each zone group is mapped to each of the host systems. SAS zoning is increasingly being targeted at users or customers that do not necessarily have storage area network knowledge or skills and thus, do not have the necessary background to perform the manual zoning of SAS based storage area networks. Thus, these users or customers require tools to help make the zoning of the storage area network easier.\nHowever, SAS poses some unique zoning challenges that are not present in more conventional Fibre Channel architecture based SANs. For example, in Fibre Channel, defining a zone configuration is simply a matter of grouping ports into a “zone” such that all ports in the “zone” can see each other. SAS zoning, on the other hand, is a two-step process. It requires the user to group ports into a zone group but then a separate step is required for access. Another complexity is the limit on the number of zone groups that can be defined. The current limit is 128 but is moving up to 256. This means that certain techniques like assigning each initiator port to its own zone (a technique common in Fibre Channel zoning) will not work in larger SAS SANs. Another example of complexity is the manner in which SAS zoning boundaries (ZSPDS) are defined. For SAS, access to devices beyond the SAS domain zoning boundary is defined and allowed. Thus, it is a complex task to understand and comply with the rules for access in SAS SANs. For Fibre Channel, the zone graph ends at the zoning boundary, i.e. there is absolutely no access beyond that boundary.\nAnother emerging aspect of SAS that makes it more complex is that it resides inside and outside the server or blade system chassis. The SAS topology is not simply a graph of external nodes as in Fibre Channel. It is more hierarchical, with a mixture of internal and external nodes. Thus, any SAS zoning mechanism will have to accommodate the different rules for the portions of the SAS SAN within the server or blade system chassis and the portions that are external to the server or blade system chassis. Thus, Fibre Channel zoning mechanisms typically make use of a single, monolithic graph for managing the Fibre Channel based SAN. Such an approach cannot be applied to SAS SANs due to the various complexities mentioned above.\nSAS zoning manages the physical connections hosts make to storage. Access management is controlled through logical unit number (LUN) mapping and masking. A logical unit number is an address for an individual logical disk drive. The term LUN is used in the SCSI protocol as a way to differentiate individual logical disk drives within a common SCSI target device like a disk array. The term has become common in storage area networks (SAN). Today, LUNs are normally not entire disk drives but rather virtual volumes. Thus, a particular storage subsystem may have hundreds, or perhaps more, LUNs. Thus, storage may have several levels of granularity from the chassis level to the blade level to the storage subsystem level to the disk drive level to the logical volume level.\nLogical unit number masking is an authorization process that makes a LUN available to some hosts and unavailable to other hosts. LUN masking is mainly implemented not as a security measure per se, but rather as protection against misbehaving servers from corrupting disks belonging to other servers. For example, Windows® servers attached to a SAN may under some conditions corrupt non-Windows® volumes on the SAN by attempting to write Windows® volume labels to them. By hiding the other LUNs from the Windows® server, this can be prevented, since the Windows® server does not even realize the other LUNs exist. “WINDOWS” is a trademark of Microsoft Corporation in the United States, other countries, or both.\nThus, an administrator may configure a SAN with two distinct steps: zoning for the physical connections between hosts and storage and LUN masking/mapping for access control at the storage manager."}
{"text": "1. Field of the Invention\nThe present invention relates to an ink supply device for a rotary printer or a sheet-feed printer and an ink key thereof, and in particular, to an ink supply device provided with an ink tray removably mounted in an ink box and an ink tray thereof.\n2. Description of Related Art\nA printer such as a sheet-feed printer or a rotary printer, as shown in FIG. 25 and FIG. 26, is provided with an ink supply device 60 with an ink box (ink fountain) 62 before a primary ink roller (ink fountain roller) 61. The ink box 62 is constituted by the peripheral surface of the primary ink roller 61, a plurality of ink keys (blades) 63 forming the bottom portion of the ink box 62, and two side plates 64 (only the deep side plate is shown in FIG. 25) which are mounted on both outer sides of the outermost ink keys of the plurality of ink keys 63 and whose front ends are in sliding contact with the peripheral surface of the primary ink roller 61. Ink in the ink box 62 is supplied to the primary ink roller 61 from a gap between the primary ink roller 61 and the tip ends of the ink keys 63 and is transferred to a group of ink rollers arranged downstream via a drawing roller not shown.\nAs shown in FIG. 26, the plurality of ink keys 63 are arranged in parallel in the direction of width of the device and the neighboring ink keys 63, 63 are in sliding contact with each other, and the side end ink keys 63 at both side ends and the side plates 64 are also in sliding contact with each other. Further, each ink key 63 can be oscillated individually around a fulcrum shaft 65 and an ink quantity controller 66 is mounted below each ink key 63.\nThe ink quantity controller 66 is provided with a push-up member 66b engaging with the bottom surface of the tip end of each ink key 63 and a pusher 66a contacting the push-up member 66b and extending or contracting to oscillate the push-up member 66b. The push-up member 66b is oscillated by extending or contracting the pusher 66a to move up and down a portion engaging with the ink key 63, whereby the tip end of the ink key 63 is oscillated. The gap between the ink key 63 and the primary ink roller 61 is controlled by the oscillation to control the thickness of an ink film supplied to the primary ink roller 61.\nFIG. 27 to FIG. 30 schematically show the structure of a conventional ink storage device of a printer mounted in a sheet-feed printer. FIG. 27 shows a state of operation and FIG. 28 shows a state of cleaning and FIG. 29 is a partial perspective view and FIG. 30 is a side view of the ink keys when they are cleaned.\nIn each drawing described above, reference numeral 101 designates an ink key controlling the amount of ink supplied and a plurality of ink keys are arranged in the direction of axis of the primary ink roller 102, the number of the ink keys being determined by the necessity of controlling the amount of ink in the direction of width of a printed matter. Reference numeral 109 designates a turning fulcrum shaft of the ink key 101 when the ink key 101 is controlled. Reference numeral 102 designates the primary ink roller for receiving the controlled amount of ink and transferring the ink to the next roller. Reference numeral 124 designates a gap formed between the ink key 101 and the primary ink roller 102 for controlling the amount of ink to be supplied. Reference numeral 111 designates ink box side plates arranged on opposite ends of the primary ink roller 102. Each ink box side plate 111 contacts the surface of each end of the primary ink roller 102 at the tip end thereof and the side surface of the ink key 101 arranged at right and left side ends at the side surface thereof to prevent the leakage of ink from these contact portions. This way, the ink key 101, the primary ink roller 102 and the ink box side plate 111 constitute an ink box 100 storing the ink.\nAn ink key receiving base 108 supports the ink key 101 and the ink box side plate 111 and is supported by a turning center shaft 110 mounted on a mechanical frame and described below. Reference numeral 107 designates a mounting bolt arranged in a groove 108a made in the ink key receiving base 108 and screwed into the bottom surface of the ink key 101. Reference numeral 106 designates a compression spring arranged in the groove 108a made in the ink key receiving base 108 and between the ink key receiving base 108 and the mounting bolt 107. The compression spring 106 applies with the mounting bolt 107 a pressing force pressing the ink key 101 toward the ink key receiving base 108. Reference numeral 103 designates an ink quantity controller mounted on each ink key 101. When the amount of ink supplied to the primary ink roller 102 is reduced (a gap 124 is reduced), a push-up portion 104 is moved up to push up the ink key 101 against the force of the compression spring 106. When the amount of ink supplied to the primary ink roller 102 is increased (i.e., a gap 124 is increased), a push-up portion 104 moves downward to push down the ink key 101 by the force of the compression spring 106.\nThe turning center shaft 110 supports the right and left ends of the ink key receiving base 108 and acts as a turning center for separating the ink key 101 and the ink box side plate 111 backward from the primary ink roller 102, as shown in FIG. 28, when the ink in the ink box 100 is removed and the ink keys 101 and the like are cleaned. A plurality of ink keys 101 are arranged in the direction of axis of the primary ink roller 102, as shown in FIG. 29, and there is provided between the ink keys 101 a small gap allowing the individual ink keys 101 to slide.\nThe conventional ink supply device 60 shown in FIG. 25 and FIG. 26 has a small gap between the neighboring ink keys 63, 63 and a small gap between the side end ink key 63 and the side plate 64, whereby the ink keys 63 can slide. Therefore, the ink may possibly get into the small gap between the ink keys 63, 63 because of capillary phenomenon or the like. The conventional ink supply device 60 has a problem that if the ink which has entered into the gap between the ink keys 63, 63 solidifies, the ink makes the action of the ink keys 63 unstable or fixes the ink keys 63 in the worst case to make it impossible to control the thickness of an ink film with high accuracy.\nFurther, it is necessary to wipe the ink remaining in the ink box 62 with textile waste or to wash it with cleaning liquid, but it is difficult to remove the ink because the ink has high viscosity. In particular, it is difficult to remove the ink from the gap between the ink keys 63, 63, and lead to increased workload on workers cleaning the ink keys 63. Further, in order to improve productivity, it is required that a preparation time for order changes be shortened to increase the availability of the device, but a cleaning time is increased because the load of cleaning is increased when the ink is changed. Therefore, it has been required that workload be reduced in cleaning operations and that a cleaning time be shortened to increase the availability and productivity of the device.\nFurther, the ink key 101 is erected approximately 90 degrees with respect to its original position as shown in FIG. 30 and the sides thereof are cleaned. A press-down unit 105 for pressing down the ink key 101 (which is constituted by a compression spring 106, a mounting bolt 107 and the like) is required to be disassembled. However, since the printer has a great number of the press-down units 105, disassembling of the press-down units 105 becomes heavily burdensome.\nFurther, after a daily printing work finishes, the sides of the ink key 101 are cleaned by picking up the tip end of each ink key 101 with fingers without disassembling the press down units 105. However, since the ink key 101 is not completely picked up unlike FIG. 30, the sides of the ink key 101 cannot be cleaned sufficiently. Further, since the ink key 101 is picked up against the spring force of the compression spring 106, there is produced a problem that the cleaning work is burdensome.\nThe present invention has been achieved in consideration of the above described problems. It is an object of the present invention to provide an ink supply device which can prevent ink from getting into a gap between ink keys to make the action of the ink keys stable and save labor in cleaning of the ink keys, and the ink key therefor.\nFurther, it is another object of the present invention to provide an ink storage device for a printer in which the sides of the ink key are cleaned easily and sufficiently.\nIn order to accomplish the objects described above, in accordance with the one aspect of the present invention, there is provided the first aspect of an ink supply device comprising an ink box whose bottom portion is formed of a plurality of ink keys arranged in parallel to each other and whose side walls are formed of side plates arranged on opposite outer sides of the plurality of ink keys, and for supplying ink from the ink box to a primary ink roller, wherein the ink supply device further comprises an ink impervious cover member covering the surfaces of the plurality ink keys, and an inset is made into either one or both of the opposing side walls of each of the plurality of ink keys abutting on each other from the top surface of the ink key to the bottom surface thereof.\nThe second aspect of an ink supply device in accordance with the present invention is characterized in that, in the first aspect of the ink supply device, the cover member is extended to cover the side plates and an insert area is provided in one or both of the sides of the side wall of the side plate, and the ink key abutting on the side plate from the top surface of the ink key to the bottom surface thereof.\nThe third aspect of the ink supply device in accordance with the present invention is characterized in that, in the first aspect or the second aspect of the ink supply device in accordance with the present invention, the inset wall portion is positioned at the tip end portion of the ink key.\nThe fourth aspect of an ink key in accordance with the present invention is characterized in that, in the ink key forming the bottom portion of an ink box, an inset is provided in the side wall of the ink key and the inset extends from the top surface of the ink key to the bottom surface thereof.\nThe fifth aspect of an ink key in accordance with the present invention is characterized in that, in the ink key of the fourth aspect, the ink key includes a covered portion which is covered by a cover member and is not in direct contact with ink, and an exposed portion which is projected upwardly via a step nearer to the tip end than the covered portion and hence is not covered by the cover member and the step has a top surface portion in direct contact with the ink, and wherein the inset wall portion is formed nearer to the base end of the key than to the step.\nThe sixth aspect of an ink key in accordance with the present invention is characterized in that, in the ink key of the fifth aspect, the end portion of the tip end of the wall inset is made at a position where the step is formed.\nIn order to solve the problems described above, the seventh aspect of the present invention is characterized in that, in an ink storage device of a printer comprising a plurality of ink keys constituting the bottom surface of an ink box, each ink key is freely turned via a turning fulcrum shaft and receives the action of a spring force in the predetermined turning direction and is provided with an aspect for applying the spring force to the ink key or removing the spring force applied to the ink key.\nThe eighth aspect of the present invention is characterized in that, in the seventh aspect of the present invention described above, the device is provided with a spring force transmission member and the ink key has a hole allowing the head portion of the spring force transmission member to pass therethrough and preventing the head portion from passing therethrough when the head portion is turned a predetermined amount, wherein the head portion of the spring force transmission member is turned to a position where the head portion can not pass through the hole to engage the spring force transmission member with the ink key, thereby applying the spring force to the ink key, and wherein the head portion of the spring force transmission member is turned to a position where the head portion can pass through the hole to disengage the spring force transmission member from the ink key, thereby removing the spring force applied to the ink key.\nFurther, the ninth aspect in accordance with the present invention is characterized in that, in the seventh aspect of the present invention described above, the ink key has a groove allowing the head portion of the spring force transmission member to pass therethrough or preventing the head portion of the spring force transmission member from passing therethrough, depending on the rotational position of the spring force transmission member, wherein the spring force transmission member is rotated to a position where the head portion of the spring force transmission member can not pass through the groove to engage the spring force transmission member with the ink key, thereby applying the spring force to the ink key, and wherein the spring force transmission member is rotated to a position where the head portion of the spring force transmission member can pass through the groove to disengage the spring force transmission member from the ink key, thereby removing the spring forceapplied to the ink key.\nThe tenth aspect in accordance with the present invention is characterized in that, in any one of the seventh to ninth aspect of the present invention described above, an ink tray is mounted above the ink key so that it may cover at least the head portion of the spring force transmission member."}
{"text": "1. Field of the Invention\nThe present invention relates to an automatic document conveying device that conveys a document to a reading position automatically, and an image forming apparatus including the automatic document conveying device.\n2. Description of Related Art\nAs an automatic document conveying device used in a copier, a scanner, a facsimile, etc., there are well-known automatic document conveying devices in a form capable of inverting front and back surfaces of a document so as to read both surfaces of the document. In recent years, along with the downsizing of an image forming apparatus, there is also a demand for the downsizing of the automatic document conveying device as well as an image forming portion of the image forming apparatus.\nFIG. 13 illustrates a related technology of such an automatic document conveying device. The related technology of the automatic document conveying device includes a document stacking tray (101), a pickup roller (102), a sheet feeding roller (103), a conveying drum (104), a front reading roller (105) and a back reading roller (106) disposed so as to be in contact with the conveying drum (104), a branching member (107), a lower inversion roller (108), an upper discharge roller (109), an inversion guide (110), and a discharge tray (111) in this order in a conveying direction (direction of an arrow “a”) of a document (D). Further, a document reading position (R) is located below the conveying drum (104) between the front reading roller (105) and the back reading roller (106). For convenience of the description, an upper surface of the document placed on the document stacking tray (101) is referred to as an A-surface and a lower surface thereof is referred to as a B-surface.\nWhen only the A-surface of the document is read, the A-surface of the document (D) supplied by the pickup roller (102) from the document stacking tray (101) is read at the document reading position (R). The branching member (107) is supported so as to be swingable about an axis at its back end. When only the A-surface of the document (D) is read, the branching member (107) is placed at a position indicated by a solid line illustrated in the figure, and the document (D) is discharged to the discharge tray (111).\nWhen both the surfaces (A-surface and B-surface) of the document are read, after the A-surface of the document (D) conveyed from the document stacking tray (101) is read at the document reading position (R), the document (D) is inverted by the lower inversion roller (108), the upper discharge roller (109), and the inversion guide (110) due to the presence of the branching member (107) at a position indicated by a dotted line illustrated in the figure. Then, the document (D) is conveyed to the document reading position (R) and the B-surface is read. The document (D) is further inverted again due to the similar operation and discharged to the discharge tray (111).\nIn the above-mentioned related technology, the number of members such as a conveying roller for use in the automatic document conveying device is reduced by configuring the automatic document conveying device as described above. However, drive sources of the branching member (107) and the lower inversion roller (108) are provided separately, and hence there is a problem in that spaces are required for setting the respective drive sources. Further, in a case of using an inexpensive solenoid as a drive source, the noise from the drive source increases, which requires countermeasures against the noise. Consequently, there is also a problem in that the number of parts increases to increase the number of assembly steps, leading to an increase in production cost."}
{"text": "This invention relates to a high-resolution, cryogenic, side-entry type specimen stage for a high-voltage electron microscope in which in-situ four point D.C. electrical resistivity measurements and electron optical observation can be made.\nStudies of the effect of irradiation on metals are of great importance today. For example, extensive work is being done on electron damage studies of metals which might be useful in fusion reactors. For this purpose optical studies using electron microscopes are of importance. Likewise, resistivity studies are of great importance because defect damage increases the electrical resistivity which can be detected before the defects become pronounced enough to see visually. While both types of studies have been made in the past, no instruments are available within which electron optical studies and electrical resistivity measurements can be made simultaneously at cryogenic temperatures. The electrical resistivity measurements would complement the electron optical observation in the high-voltage electron microscope to yield a unique opportunity to investigate defect production in metals by electron irradiation over a wide range of defect concentration. The specimen stage described herein is designed to achieve the same electron optical resolution (.about.5 A) below 10 K as can be obtained at room temperature."}
{"text": "Many organizations poll their customers, suppliers, contractors, and employees in an attempt to determine how well the organization might be performing from the viewpoint of these entities. For instance, many merchants poll their customer base in an attempt to determine how well the organization fills orders, meets delivery deadlines, competes in terms of price, etc. Similarly, organizations some times poll their suppliers and contractors in an attempt to determine how clear their instructions are, how streamlined their procedures are, how adequate their payment methods are, etc. Likewise, employers often poll their workforce in an attempt to determine how satisfied the employees are, how well the human resource policies are implemented, and the like.\nIn many of these cases, the polling is done in a continuous manner. To conduct a continuous poll (such as a customer satisfaction survey or an employee “360” review) an organization sends out a survey to a number of potential respondents at a pre-determined frequency. Commonly, the organization conducts the poll on a monthly or quarterly basis. After sending the survey out, the organization waits for some time for the responses to come in. Eventually, the organization closes the survey and compiles the results."}
{"text": "The present invention relates to medical imaging of the heart, and more particularly, to automatic detection of anatomic landmarks of the left ventricle in magnetic resonance long axis image slices.\nCardiovascular disease is the leading cause of death in developed countries. Early diagnosis can be effective in reducing the mortality of cardiovascular disease. Quantification of the left ventricle (LV) is of particular interest among the four heart chambers because it pumps oxygenated blood from the heart to the rest of the body. In particular, precise measurements of both the dimensions and functions of the LV are essential in clinical applications for diagnosis, prognostic, and therapeutic decisions. Magnetic resonance (MR) imaging can accurately depict cardiac structure, function, perfusion, and myocardial viability, and precise measurements can be achieved using MR imaging. Accordingly, MRI is widely accepted as the standard for heart chamber quantification. However, due to the considerable amount of information available, analysis, such as segmentation, of cardiac images for functionality quantification is time consuming and error-prone for human operators. Thus, automated methods for analyzing MR images are needed.\nIn MR scans, long axis slices are used as scout images for acquisition planning, as well as to complement a stack of short axis slices. Long axis slices capture the LV's shape information and can also be used to correct mis-registration of the short axis stack. Long axis acquisitions can be an image sequence of long axis slices or a single slice that is scanned during MR acquisition planning. Anatomic landmarks in an MR long axis slice can be used for higher level segmentation, such as initialization of deformable model based approaches, and for accelerating acquisition time of a full MR scan by facilitating fully automatic planning of cardiac MR examinations. Thus, automatic detection of anatomic landmarks in a cardiac MR long axis slice is desirable."}
{"text": "The patent is directed to a multi-player game, and more particularly to a multi-player game having an electronic instruction unit providing game instructions to the players during the course of game play.\nVarious games having electronic or electro-mechanical apparatus associated therewith have been previously described. For example, U.S. Pat. No. 4,572,513 to Evans discloses an educational game including a game board, at least one game piece and an audio recording. The game board has a serpentine path located on the top face thereof. The serpentine path is broken up into a series of connected squares. The game piece is designed such that it is moved along the path as the game is played. The audio recording contains a number of mathematical problems or other questions and the answers thereto. Each question or problem is separated from its answer by a pause of predetermined duration. Each time the player of the game answers a problem or a question correctly and within the time allowed, the player advances his/her game piece along the serpentine path. The number of squares on the path and the number of questions on the recording are equal so that if a player answers all of the questions correctly, his/her game piece exactly reaches the finish space.\nU.S. Pat. No. 5,011,156 to LaChance, Jr. et al. discloses a game apparatus consisting essentially of a game board, cards, set of playing pieces for each player, a compact disc with multiple tracks set up in such a way that the playing pieces, or tokens, for each player will move around the board surface which may comprise a single or dual segmented path. Part or all of the segments contain one of the numbers in a two number selection code on the compact disc. Cards which are provided as part of the game apparatus will contain a second number of the two number selection code so that as playing the game, once both numbers are determined, the player will know which track to punch in on the compact disc control. The instructions on the compact disc track will determine further moves and actions and enhances the game play."}
{"text": "The invention relates to sequencing the operational states of a turbine, and particularly to developing sequencing algorithms for controlling the operational states for industrial gas and steam turbines.\nIndustrial gas and steam turbines typically operate at predefined operational states. With respect to a gas turbine, these states relate to starting the gas turbine, accelerating the gas turbine to a rotational speed (load speed) for driving a load for power generation or a mechanical device, and shutting down the gas turbine. As an example, the operational states during startup may include starting auxiliary devices for the gas turbine, mechanically rotating the shaft of the gas turbine, and initiating ignition of combustion in the gas turbine. Other operating states are associated with accelerating the gas turbine to an idle or no-load speed, running the gas turbine at speed and under load, and shutting down the gas turbine. An industrial gas turbine operates at one of its predefined operational states. A steam turbine will also have predefined operating states and will transition between its operating states.\nA software program generally referred to as a sequencer determines the current operational state of a turbine, selects the next operational state and determines when to transition from the current to the next operational state. The sequencer software module is conventionally stored and executed by a computer control system for the gas turbine.\nThe control system may execute other software modules which generate control commands to operate specific components of the gas turbines, typically referred to as auxiliary systems. These auxiliary systems may control the: fuel valves that regulate fuel flow to the gas turbine, starter motor that mechanically turns the compressor and turbine, instruments and sensors monitoring the gas turbine, mechanical actuators for the inlet guide vanes (IGVs), and pumps for oil and fuel. The sequencer software program communicates with the other software modules to monitor the operation of the gas turbine and notify the other modules as to the state of the turbine.\nCurrent day sequencing software modules are written for a specific turbine model or family of models. Once written, the software instructions are tested to confirm that they properly control the turbine and are free of error. The tested software instructions are documented with comments in the software coding and manuals for using and configuring the sequencing module to a specific gas turbine. The writing, testing and documentation of the sequencing module for a new model or family of models of turbines are time consuming, expensive and require software programmers and engineers familiar with the operation of the specific gas turbine.\nHistorically, sequencing modules are designed and developed specifically for each model of an industrial gas turbine. Some individual manufacturers of industrial gas turbines have developed multi-model sequencers for a related group of models, such as a product family of gas turbines. These sequencing modules are limited to the gas turbine model(s) for which they are designed. Traditionally, sequencing modules have not been adapted to control gas turbines beyond those models for which the sequencing module was initially designed.\nWhile some existing sequencing modules created for one turbine have been adapted to work on other models of turbines, this ad hoc approach to adapting sequencing modules introduces risks that the adaptation of the sequencing module does not properly sequence the new turbine through its operating states. This ad hoc approach is not an efficient approach for developing sequencers for a large number of turbine models.\nBecause they are custom developed for each new model or family of models of gas turbines, the sequencing modules for different models/families of gas turbines have large variations in their software structure and software instructions. To work with the sequencing modules of different models/families of gas turbine requires knowledge of the software in each module. A person qualified to work on a sequencer for one gas turbine model may not be qualified to work on the sequencer of another gas turbine module or may not be knowledgeable of subtle but important differences between sequencers for different turbines.\nIn view of the cost and time required to develop a gas turbine sequencer and the variations between sequencers for different models, there is a long felt and unmet need for methods and systems to reduce the cost and time required to develop sequencing modules for a wide range of models of a gas turbine."}
{"text": "The present invention relates generally to the field of digital image synthesis and, more particularly to digital color image reproduction, digital color image printing, digital color image display and anti-counterfeiting protection. The invented method and apparatus enable generating digital halftone images by multi-color dithering. The invented method and apparatus are particularly effective for generating images on color printers printing with more colors than cyan, magenta, yellow and black, and for generating images incorporating microstucture patterns (micro-characters, symbols and ornaments) made of a combination of standard and non-standard inks. Non-standard inks comprise non-process colors, opaque inks, and special inks such as metallic ink, variable color inks and fluorescent inks. The generated images are displayed on color displays or printed on physical output devices such as ink-jet printers, electrophotographic printers, offset printers or gravure printers.\nThe reproduction of color images requires in the general case (1) separating the image colors (for example red, green and blue) into the set of available printable colors (for example cyan, magenta, yellow and black), (2) halftoning each of the printable color layers and (3) printing the halftoned image or visualizing it on a display device.\nMost existing color halftoning techniques are based on the independent halftoning of each of the contributing color layers. In offset and in many electrographic printers, the color halftone layers are generated independently at angles of 15, 45, 75 degrees in order to avoid interferences between the cyan, magenta and black layers. Ink-jet printers often use error-diffusion to halftone each of the color layers independently.\nExisting approaches for color separation with non-standard inks consider that color layers can be printed independently from one another. This is true as long as the inks are transparent and as long as the dots of each screen are randomly positioned with respect to the dots of the other screens, as assumed by the Neugebauer equations (see G. L. Rogers, “Neugebauer Revisited: Random Dots in Halftone Screening”, Color Research and Applications, Vol. 23, No. 2, pages 104–113, 1998).\nIf these assumptions are not true, halftoning each of the color layers independently may generate color shifts depending on the amount of superposition between screen dots of individual halftone layers. As long as the inks are transparent, the color shifts are small and the reproduction fidelity can be ensured by calibrating the printing device (see H. R. Kang, Color Technology for Electronic Imaging Devices, SPIE Publication, 1997, chapter 12).\nSome applications however require that different inks be always printed side by side without overlapping. This is the case, for example, when printing with opaque inks as taught by inventor H. Kueppers in U.S. Pat. No. 4,812,899, issued Mar. 14, 1989. Different inks must also be printed side by side for some printing devices as taught by J. Gereadts and S. Lenczowski, “Océ's productive colour solution based on the Direct Imaging Technology”, Proceedings of the IS&T International Conf. On Direct Imaging Technology, Vol. NIP-13, 1997, pages 728–733.\nManufacturers of valuable documents, such as banknotes, identity cards and checks make often use of the high registration accuracy of their original printing equipment to create a graphic design which is both visually pleasant and difficult to imitate using standard printing processes (see Rudolf L. VanRenesse, Optical Document Security, Artech House 1998, chapter 6). In addition, they may use non-standard inks (sometimes also called custom inks), i.e. inks which have a color different from process colors Cyan, Magenta, Yellow and Black to reproduce images and graphic designs. Similar anti-counterfeiting techniques are used to print commercial packages which need to be protected against illegal counterfeits (drugs, cosmetics, etc.).\nOne prior art method of generating a dither pattern for an RGB value, is taught by D. N. Weise and H. Gunter Zieber in U.S. Pat. No. 5,485,558, issued Jan. 16, 1996. That method is limited to applications, mainly display applications, where one can associate to each input color pixel a dither array, for example an 8×8 array, in order to reproduce that color pixel. Since the dither array must remain of small size, that method does not allow to create at the screening layer artistic screen shapes. Another prior art method of color halftoning is to use vector error-diffusion in color space, as taught by J. Sullivan and R. L. Miller in U.S. Pat. No. 5,070,413, issued Dec. 3, 1991. However, error-diffusion has the drawback of inducing considerable dot gain, compared with clustered dot dithering methods. In addition, error-diffusion does not allow to create user-defined screen shapes.\nThe present invention discloses a new multi-color dithering method which automatically enforces side by side printing of several color and/or metallic ink layers. Thanks to a user definable dither function or dither array of freely chosen size, multi-color dithered images may incorporate identifyable screen shapes, for example artistic screen shapes, which can be used for document authentication. Combined with color separation by tetrahedral interpolation, multi-color dithering enables printing images made of non-standard inks.\nThe invented method is a non-trivial generalization of ordered dither halftoning, which uses spatial dithering in order to obtain clustered-dot or dispersed-dot binary patterns on output devices (see H. R. Kang, Digital Color Halftoning, SPIE Press/IEEE Press, 1999, chapters 13 and 14)."}
{"text": "There is a growing concern about carbon source after the depletion of fossil fuels in the world of polymeric materials, as the majority of the commercialized polymer products are related to petrochemical resources, such as polystyrene, polyvinyl chloride, polypropylene, and polyethylene terephthalate. Renewable resources, such as carbohydrate, lignin, cellulose and plant oils, gain attention for the production of polymeric materials, due to their short term of reproducibility from biomass. Plant oils, mainly constituting of triglycerides, are extracted primarily from the seeds of oilseed plants. Besides the basic consumption as food, other applications are developed from this bio-renewable resource, like lubricants, biofuels, plasticizers, and construction materials. These triglycerides are also candidates as the building block of polymeric materials considering their natural abundance and inherent functionality, like the unsaturated double bond and ester groups. Challenges still present, including their heterogeneous and non-uniform structures.\nVarious kinds of biomass have been applied to prepare polymeric materials, such as cellulose, hemicellulose, lignin, starch, chitin, rosin acid, lactic acid, and plant oils. Among them, plant oils are one of the most abundant raw materials, which can be obtained primarily from plant seeds. Due to the unsaturation of fatty acid chains, plant oils have been widely used in the preparation of thermoset polymeric materials. Cationic polymerization of conjugate low saturation plant oils or copolymerization with divinylbenzene leads to thermosetting polymers. However, the reactivity of the unsaturated fatty acid double bonds is not high enough to obtain strong thermosetting polymers.\nAnother important method is to transfer the double bonds of plant oils to epoxy groups. The epoxy group can be further modified to obtain high reactive multi-functional monomers, such as acrylate plant oils. These acrylate plant oils can be used in the design of thermoset polymers. Epoxy resins are unique among all thermosetting resins with excellent resistance to chemicals and good thermal properties, which are widely used in coatings, adhesives, and matrices for composite materials. Epoxidized soybean oils are commercially available and have been used as plasticizers in plastic industry. However, the preparation and application of epoxy resins from plant oils is still limited, due to the poor mechanical properties of the epoxy resins derived from plant oils. There are two main reasons: the reactivity of the plant oil epoxy groups is too low; the number of epoxy groups in an epoxidized plant oil is too low."}
{"text": "1. Field of the Invention\nThe invention relates generally to the field of social web sites on global computer networks and in particular to a system and method for managing information personal to individual users of social web sites.\n2. Statement of the Problem\nThe Internet started as an information source and communication tool, and these functions remain important aspects of the Internet. It is also a truism that the Internet is a powerful sales tool. Companies that provide useful internet services have cash flows of more then ten billion dollars a year, much of which comes from advertising. The Internet has become the most important medium for distribution of music, and is rapidly becoming the most important medium for distribution of videos. The Internet has also become an important vehicle for telecommunications, including telephone communication. Because the Internet lends itself to immediate feedback to advertisers and service providers, new Internet uses, such as directed advertising are rapidly evolving.\nThe Internet has also grown to become an important social network. From one point of view, the Internet is the largest social network on the planet. Moreover, on a more explicit level, there are scores, if not hundreds of social network sites on the Internet, and millions of blogs. Social network sites are defined as web sites or applications that provide people with designated relationships to each other access to information that is not available to others who do not have such designated relationships. Because the Internet can be accessed from a myriad of points, such as desktop computers connected to wires, laptop computers connected to wireless networks, hand-helds, cell phones, and integrated products such as Ipods™, for many individuals, the Internet has become their primary means of communication. The explicit social web sites encourage individuals to share their ideas, experiences and thoughts with their friends, sometimes on an hourly or more frequent basis. These social web sites are uniquely different from conventional informational and shopping web sites in that they provide interactions between the web site users, and not just between the web site and individual users. The actions on these web sites are more complex, including actions such as adding a contact, deleting a contact, sending a message to a contact, inviting, such as to a social event, blocking, staring, which also is called bookmarking or placemarking, depending on the web site, linking as in assigning an attribute or giving credit, set a trust level, and voting.\nThe Internet permits individuals to live in many different social circles that would not be available without the Internet. However, it also creates problems in that it is difficult to control identities on the internet. A person may select an identity, and find that the identity may be taken by another.\nAn even greater problem is created by the fact that after an individual or company registers an identity on one site, others may intentionally usurp that identity on another site. The proliferation of sites makes it difficult for an individual or company to capture an identity on all sites.\nAll of the above problems are particularly exacerbated in social network sites. The key purpose of such sites is social networking, and if a user's identity is not known or varies, such social networking becomes difficult to impossible. Moreover, reconnecting among friends on new social network sites is laborious and thus seriously impedes social networking.\nFor the above and other reasons, it would be highly desirable if a method and apparatus were available that allowed an individual, corporation or other entity to better control their identities on the internet."}
{"text": "This application is based on application No. 2000-47471 filed in Japan, the content of which is hereby incorporated by reference.\n1. Field of the Invention\nThe present invention relates to image data editing.\n2. Description of the Prior Art\nIn performing taking photographs with a camera, there are cases where the camera inclines due to a shake of the camera caused by a hand movement or the like. The inclination of the camera causes a problem when a subject including a horizontal straight line such as the horizon or a vertical straight line such as a tall tree or a high-rise building is photographed. That is, when a horizontal straight line or a vertical straight line is photographed with the camera inclined, a disagreeable print is obtained in which the subject is inclined with respect to the horizontal or the vertical direction of the photographic frame. In the case of a digital image taking with a digital camera or the like, a disagreeable image is obtained in which the horizontal or the vertical straight line is inclined with respect to the display frame of the monitor.\nSuch an image can be corrected so that the horizontal or the vertical straight line in the image is parallel to the display frame of the monitor by capturing the image in a computer and rotating it by use of an image editing software program. Specifically, the user captures the digital image into a computer, and enters a rotation angle into the image editing software program (for example, xe2x80x9cPhotoshopxe2x80x9d by Adobe Systems Incorporated). Then, the image is displayed being rotated by the rotation angle, so that the horizontal or the vertical straight line is parallel to the horizontal direction or the vertical direction of the display frame of the monitor.\nFor the user to enter the rotation angle of the image, it is necessary to determine how much the subject is inclined. However, it is not easy to determine the angle at which the image displayed on the monitor is inclined. It is necessary for the user to first enter a rotation angle considered appropriate and viewing the rotated image, further enter another rotation angle to correct the inclination. It is cumbersome for the user to find out the rotation angle through such trial and error.\nOther methods of determining how much the subject is inclined include determining the subject inclination by the user specifying two points on a horizontal or a vertical straight line to identify the horizontal or the vertical straight line and calculating the inclination of the straight line. To correct the subject inclination, the image is rotated so that the straight line is parallel to the horizontal or the vertical direction of the display frame of the monitor.\nHowever, even by this method, the subject inclination cannot accurately be determined when the straight line determined by the specified two points and the horizontal or the vertical straight line do not coincide with each other, that is, when the specified two points are not situated on the horizontal or the vertical straight line. Therefore, when the image is rotated based on the obtained inclination, the horizontal or the vertical straight line in the rotated image is not parallel to the horizontal or the vertical direction of the display frame of the monitor.\nAn object of the present invention is to provide a device and a method for detecting the inclination of a subject.\nA first aspect of the present invention provides an inclination detecting device comprising:\na display for displaying on a screen a subject image comprising a plurality of pixels;\na pixel specifier for a user to specify a first pixel from among the pixels;\na calculator for defining a pixel group constituting a circumference of an area including the specified first pixel, and calculating deviations in pixel value between adjoining pixels of the pixel group; and\na detector for selecting a second pixel and a third pixel based on the calculated deviations, and detecting an inclination of a straight part of the subject based on the selected second and third pixels. Moreover, the first aspect of the present invention provides an image processing apparatus provided with the inclination detecting device.\nAccording to this structure, when the user specifies the first pixel, an area including the first pixel is defined, and based on the pixel value deviations of the pixel group surrounding the area, the second pixel and the third pixel are selected. Based on the selected second and third pixels, more concretely, the inclination of the straight line determined by the second and the third pixels is detected as the inclination of the straight part of the subject. Since the inclination of the straight part of the subject represents the overall inclination of the subject in the image, how much the subject is inclined can easily be detected.\nThe detector selects pixels where the deviations are maximum and minimum as the second pixel and the third pixel.\nThe detector detects the inclination of a straight line connecting the second pixel and the third pixel as the inclination of the straight part of the subject, and the inclination detecting device further comprises an image rotation processor for rotating the image based on the detected inclination of the straight part to thereby correct the inclination.\nAccording to this structure, the image is rotated based on the detected inclination of the subject, whereby the inclination of the subject can be corrected. Consequently, only by specifying the first pixel, the user can obtain an image where the inclination is corrected.\nThe image rotation processor decides the direction in which the image is rotated, according to the inclination of the straight part.\nAccording to this structure, the direction in which the image is rotated is decided according to the detected inclination of the subject. With this, the direction in which the image is rotated is decided and the inclination of the subject is corrected without the need for the user to specify the direction. Consequently, the load on the user is reduced.\nAnother aspect of the present invention provides an inclination detecting method comprising:\ndisplaying on a screen a subject image comprising a plurality of pixels;\naccepting a specifying, by a user, of a first pixel from among the pixels;\ndefining a pixel group constituting a circumference of an area including the specified first pixel and calculating deviations in pixel value between adjoining pixels of the pixel group;\nselecting a second pixel and a third pixel based on the calculated deviations; and\ndetecting an inclination of a straight part of the subject based on the selected second and third pixels. Moreover, this aspect of the present invention provides an image processing apparatus executing the inclination detecting method.\nAccording to this method, when the user specifies the first pixel, an area including the first pixel is defined, and based on the pixel value deviations of the pixel group surrounding the area, the second pixel and the third pixel are selected. Based on the selected second and third pixels, more concretely, the inclination of the straight line determined by the second and the third pixels is detected as the inclination of the straight part of the subject. Since the inclination of the straight part of the subject represents the overall inclination of the subject in the image, how much the subject is inclined can easily be detected.\nIn the selecting, pixels where the deviations are maximum and minimum are selected as the second pixel and the third pixel.\nIn the detecting, the inclination of a straight line connecting the second pixel and the third pixel is detected as the inclination of the straight part of the subject, and the inclination detecting method further comprises rotating the image based on the detected inclination of the straight part to thereby correct the inclination.\nAccording to this method, the image is rotated based on the detected inclination of the subject, whereby the inclination of the subject can be corrected. Consequently, only by specifying the first pixel, the user can obtain an image where the inclination is corrected.\nIn the rotating, the direction in which the image is rotated is decided according to the inclination of the straight part.\nAccording to this method, the direction in which the image is rotated is decided according to the detected inclination of the subject. With this, the direction in which the image is rotated is decided and the inclination of the subject is corrected without the need for the user to specify the direction. Consequently, the load on the user is reduced.\nStill another aspect of the present invention provides a computer program product for executing a method comprising:\ndisplaying on a screen a subject image comprising a plurality of pixels;\naccepting a specifying, by a user, of a first pixel from among the pixels;\ndefining a pixel group constituting a circumference of an area including the specified first pixel and calculating deviations in pixel value between adjoining pixels of the pixel group;\nselecting a second pixel and a third pixel based on the calculated deviations; and\ndetecting an inclination of a straight part of the subject based on the selected second and third pixels. Moreover, this aspect of the present invention provides an image processing apparatus executing such a computer program.\nAccording to this computer program product, when the user specifies the first pixel, an area including the first pixel is defined, and based on the pixel value deviations of the pixel group surrounding the area, the second pixel and the third pixel are selected. Based on the selected second and third pixels, more concretely, the inclination of the straight line determined by the second and the third pixels is detected as the inclination of the straight part of the subject. Since the inclination of the straight part of the subject represents the overall inclination of the subject in the image, how much the subject is inclined can easily be detected.\nIn the selecting, pixels where the deviations are maximum and minimum are selected as the second pixel and the third pixel.\nIn the detecting, the inclination of a straight line connecting the second pixel and the third pixel is detected as the inclination of the straight part of the subject, and the computer program product further comprises rotating the image based on the detected inclination of the straight part to thereby correct the inclination.\nAccording to this computer program product, the image is rotated based on the detected inclination of the subject, whereby the inclination of the subject can be corrected. Consequently, only by specifying the first pixel, the user can obtain an image where the inclination is corrected.\nIn the rotating, the direction in which the image is rotated is decided according to the inclination of the straight part.\nAccording to this computer program product, the direction in which the image is rotated is decided according to the detected inclination of the subject. With this, the direction in which the image is rotated is decided and the inclination of the subject is corrected without the need for the user to specify the direction. Consequently, the load on the user is reduced.\nThese and other objects, advantages and features of the invention will become apparent from the following description thereof taken in conjunction with the accompanying drawings, which illustrate specific embodiments of the invention."}
{"text": "This invention relates in general to folding cartons and, in particular, to a folding carton which is adapted to be used for containing and packaging smaller individually pre-packaged items within the container.\nMore specifically, but without restriction to the particular uses which are shown and described, this invention relates to a carton formed from an open-ended tube in which the closures thereof may be formed from an integral portion of the tube body. Portions of the packaged items contained within the carton may be observed through the carton to determine such things, for example, as size, color or \"UPC\" coding.\nIn the retail trade, it is quite common to merchandise such articles as pre-packaged hoisery, smoking tobacco, office supplies, dry goods, drug supplies and other sundry items in individual packages to be purchased by a retail customer. Such items are frequently purchased by a retailer from a wholesale distributor in smaller quantities than in large case lots containing a greater quantity of these items. Therefore, it has become desirable to provide an intermediate packing container or carton within such larger packing cases or shipping containers. Such intermediate packing containers or cartons are referred to in the trade as a \"distributor pack\" and contain a smaller quantity of the individually pre-packaged consumer or retail items.\nDistributor packs used in this manner are especially useful for packing irregular shaped individual items. Such items, for example, as blister-packed tape dispensers are difficult to package as individual items in case lot quantities due to their irregular shape. Therefore, smaller quantites of such items are frequently packaged in a container, such as a distributor pack, which is used as an intermediate packer in a larger case to facilitate orderly packing and shipping of such items in case lot quantites. Such distributor pack cartons may also be utilized to conveniently pre-package smaller quantites of individually packaged items by, for example, size or color. In this manner less than case lot quantities of a particular size or color of an item may be economically packaged, or a case lot quantity containing a mixed assortment of items may be packed eliminating the necessity of case-lot purchases of each size or color.\nIn order to perform these functions, a distributor pack must provide an economical packaging medium for distribution of moderate quantities of individual pre-packaged items in order to minimize the additional packaging cost to the product manufacturer. In addition, such packaging must be sufficiently strong to provide adequate protection for the individual pre-packaged products contained therein, and to prevent loss or damage to the individual pre-packaged items or their package.\nThe present invention provides an economical distributor pack carton which does not require the addition of paperboard closure devices to close the bottom of the carton, or require the use of two-piece telescoping cartons as heretofore necessary. The distributor pack carton of the present invention has the bottom closure thereof formed from an integral portion of the carton body and protects individual packages within the carton against damage.\nThe carton body is formed as an open end tube which may be readily constructed using high speed gluing techniques since the tube is formed by single-line strip gluing along a single straight edge. In addition, the bottom closure of the open ended tube may be quickly and conveniently set up, to enable the distributor pack to be quickly filled with individual pre-packaged items at a minimum of time and expense. The open ended tube construction having the bottom closure formed from an integral portion of the tube body eliminates the necessity of complex packaging assembled about an array of individual consumer packaged goods, or expensive and complex operations requiring on-site gluing or the like.\nIn one embodiment of the invention, predetermined portions of the panels from which the open end tube carton is formed are removed to display a portion of the individual pre-packaged items to be contained within the pack. The provision of such openings in the distributor pack carton permits the ready determination of such things, for example, as size, quantity, color, price, date or \"UPC\" coding on the individual packages. In addition, providing these openings in the body of one carton allows two such distributor pack cartons to be nested together and formed from a single standard sheet of paperboard. Since standard sheets of paperboard are customarily formed in a square or rectangular configuration, individual carton blanks may be inter-nested or laid out on the standard sheet of paperboard in complementary form to provide maximum utilization of the paperboard in forming the carton blank. The formation of a top closure of one of the open ended tubes from the paperboard material removed to form the openings in the panels of the other complementary inter-nested carton blank permits the greatest efficiency and utilization of the paperboard stock from which these cartons are formed and substantially eliminates wasted paperboard."}
{"text": "Ionizing radiation has many medical, industrial and military uses. Although ionizing radiation can be used in the therapy of diseases such as cancer, exposure to biologically significant levels of such radiation can also cause genotoxic stress. Similarly, many industrial processes (such as the production of nuclear power) and military uses (such as nuclear weapons) can expose individuals to hazardous levels of ionizing radiation. Such radiation can elicit a variety of cellular responses, ranging from cell-cycle arrest to mutation, malignant transformation, or cell death. Many of the responses (such as genotoxicity) are often subtle, and exposed persons may be unaware or unsure if they have been exposed. Moreover, it may require years to evince an untoward effect (such as the development of a malignancy) caused by the exposure.\nMany of the assumptions about low dose effects have been based on extrapolations from effects measured at high doses. Transcriptional responses to doses of ionizing radiation with relatively little effect on cell survival have not been as well investigated, although small variations in expression levels of several isolated genes have been detected. A dose of 50 cGy reportedly reduced expression of β- and γ-actin (Woloschakand Chang-Liu, Int. J. Radiat. Biol. 59:1173-83, 1991) and induced RB-1 and H4 histone (Woloschak and Chang-Liu, Cancer Lett. 97:169-75, 1995) in Syrian Hamster Embryo cells, while a decrease in c-myc and increase in c-jun was detected in these cells following a dose as low as 6 cGy (Woloschak and Chang-Liu, Cancer Lett. 97:169-75, 1995). In a transformed human lymphoblast cell line, activation of NF-κB has been reported with as little as 10 cGy of radiation (Prasad et al., Radiat. Res. 138:367-72, 1994), along with induction of c-FOS, c-JUN, c-MYC and c-Ha-RAS in the 25-200 cGy range (Prasad et al., Radiat. Res. 143:263-72, 1995). The induction by 25 cGy of PBP74, a member of the heat shock 70 gene family, has also been reported in two human cancer cell lines (Sadekova et al., Int. J. Radiat. Biol. 72:653-60, 1997)."}
{"text": "Regenerative medicine harnesses, in a clinically targeted manner, the ability of regenerative cells, e.g., stem cells and/or progenitor cells (i.e., the unspecialized master cells of the body), to renew themselves indefinitely and develop into mature specialized cells. Stem cells are found in embryos during early stages of development, in fetal tissue and in some adult organs and tissue. Embryonic stem cells (hereinafter referred to as “ESCs”) are known to become many if not all of the cell and tissue types of the body. ESCs not only contain all the genetic information of the individual but also contain the nascent capacity to become any of the 200+ cells and tissues of the body. Thus, these cells have tremendous potential for regenerative medicine. For example, ESCs can be grown into specific tissues such as heart, lung or kidney which could then be used to repair damaged and diseased organs. However, ESC derived tissues have clinical limitations. Since ESCs are necessarily derived from another individual, i.e., an embryo, there is a risk that the recipient's immune system will reject the new biological material. Although immunosuppressive drugs to prevent such rejection are available, such drugs are also known to block desirable immune responses such as those against bacterial infections and viruses.\nMoreover, the ethical debate over the source of ESCs, i.e., embryos, is well-chronicled and presents an additional and, perhaps, insurmountable obstacle for the foreseeable future.\nAdult stem cells (hereinafter interchangeably referred to as “ASCs”) represent an alternative to the use of ESCs. ASCs reside quietly in many non-embryonic tissues, presumably waiting to respond to trauma or other destructive disease processes so that they can heal the injured tissue. Notably, emerging scientific evidence indicates that each individual carries a pool of ASCs that may share with ESCs the ability to become many if not all types of cells and tissues. Thus, ASCs, like ESCs, have tremendous potential for clinical applications of regenerative medicine.\nASC populations have been shown to be present in one or more of bone marrow, skin, muscle, liver and brain. However, the frequency of ASCs in these tissues is low. For example, mesenchymal stem cell frequency in bone marrow is estimated at between 1 in 100,000 and 1 in 1,000,000 nucleated cells Thus, any proposed clinical application of ASCs from such tissues requires increasing cell number, purity, and maturity by processes of cell purification and cell culture.\nAlthough cell culture steps may provide increased cell number, purity, and maturity, they do so at a cost. This cost can include one or more of the following technical difficulties: loss of cell function due to cell aging, loss of potentially useful cell populations, delays in potential application of cells to patients, increased monetary cost, increased risk of contamination of cells with environmental microorganisms during culture, and the need for further post-culture processing to deplete culture materials contained with the harvested cells.\nMore specifically, all final cell products must conform with rigid requirements imposed by the Federal Drug Administration (FDA). The FDA requires that all final cell products must minimize “extraneous” proteins known to be capable of producing allergenic effects in human subjects as well as minimize contamination risks. Moreover, the FDA expects a minimum cell viability of 70%, and any process should consistently exceed this minimum requirement.\nWhile there are existing methods and apparatus for separating cells from unwanted dissolved culture components and a variety of apparatus currently in clinical use, such methods and apparatus suffers from a significant problem—cellular damage caused by mechanical forces applied during the separation process, exhibited, for instance, by a reduction in viability and biological function of the cells and an increase in free cellular DNA and debris. Furthermore, significant loss of cells can occur due to the inability to both transfer all the cells into the separation apparatus as well as extract all the cells from the apparatus. In addition, for mixed cell populations, these methods and apparatus can cause a shift in cell profile due to the preferential loss of larger, more fragile subpopulations.\nThus, there is a need in the field of cell therapy, such as tissue repair, tissue regeneration, and tissue engineering, for cell compositions that are ready for direct patient administration with substantially high viability and functionality, and with substantial depletion of materials that were required for culture and harvest of the cells. Furthermore, there are needs for reliable processes and devices to enable production of these compositions that are suitable for clinical implementation and large-scale commercialization of these compositions as cell therapy products."}
{"text": "Compositions used to enhance cosmetic products are known in the art. Such compositions include those that are applied over top compositions such as lipstick to provide attributes such as gloss, lubricity and transfer-resistance of the cosmetic product they are applied over. These enhancement products utilize a variety of polymeric fluids and film forming technologies. For example, acrylic film-formers that are incorporated in lipstick overcoat products such as CSI Incorporated's \"Sealed with a Kiss\" are delivered in a volatile vehicle, alcohol, which is spread over the lipstick surface.\nAlternative topcoat products to those described above are disclosed in Japanese Pat. Application Number HEI 5 [1993]-221829, published Aug. 31, 1993. Said overcoats are reputed to exhibit improved durability of makeup effect, suppression of color transfer, and improved applicability. Said topcoats comprise from 0.2 to 25% of silica powder and/or alumina powder and from 75% to 99.8% of a perfluoropolyether of general formula: ##STR1## wherein R.sup.1 though R.sup.5 are independent fluorine atoms, perfluoroalkyl groups, or oxyperfluoroalkyl groups; the value of p, q, and r is at least zero; wherein the perfluoropolyether molecular weight is from about 500 to about 10,000, wherein P, Q and R may be equal, but, not zero. The preferred perfluoropolyether disclosed therein is a commercially available product known as Fomblin HC-04, HC-25, and HC-R available from Montefluosu of Milano, Italy.\nWhile such compositions may provide certain advantages, it has been found that they often disrupt the primary advantages of the cosmetic products they are applied over. For example, cosmetic products compromise their gloss or feel attributes in order to improve the long wear properties provided by the composition that is applied over top the cosmetic product. Alternately, cosmetic products must sacrifice long wear properties in order to improve the gloss and or feel attributes provided by the such compositions."}
{"text": "A network processor is a device that executes programs to handle packet traffic in a data network. A network processor is also often referred to as a network flow processor or simply a flow processor. Examples include network processor integrated circuits on router line cards and in other network equipment. In one example, a network processor integrated circuit is capable of receiving packets, classifying and performing atomic operations on the packets and associated packet data, and transmitting packets. Processors on the integrated circuit are called upon to perform processing functions that include managing various queues of data to ensure various levels of Quality of Service (QOS) to various network clients. A processor on the integrated circuit may also be called upon to determine and to log updated activity and permission information into appropriate tables in memory. As throughput requirements increase, ways of adding processing power are sought.\nIn one specific example, a network processor integrated circuit uses the flexible and expandable IXP2800 architecture. The IXP2800 architecture allows multiple high-speed processors (referred to as microengines) to access the same memory resources via a common command/push/pull bus. Due to use of the IXP2800 architecture and multiple microengines, increased processing power is brought to bear on the tasks of identifying data structures using hash functions and of logging packet and byte count information. If more throughput is required, then more microengines can be employed. If less throughput is required, then fewer microengines can be employed. The NFP-3XXX and NFP-6XXX families of network processor integrated circuits available from Netronome Systems, Inc. of Santa Clara, Calif. include a selection of IXP2800-based network processor integrated circuits having different numbers of microengines."}
{"text": "As of now, in the quest to intensify the iron-making process the efforts are directed to improving the blast-furnace process through optimum distribution of a gas flow over a shaft section of the blast furnace, which determines the quality of iron produced, the specific fuel consumption for its production and the blast furnace output as a whole. The uniform distribution of the gas flow over the furnace section is attained by controlling the rate of blowing and changing the layer contours of ore and coke components of the burden which is fed into a throat zone by a charging device installed on the furnace cone. Therefore, the requirements that are placed upon the charging device are to provide a desired contour of the burden layer over the furnace cross section with minimum nonuniformity in the circumferential grain-size distribution of the burden in the throat zone, and, essentially, a controlled thickness of the burden layer charged.\nConventional blast furnace charging apparatuses include a double-cone charging apparatus equipped with movable throat plates, which is disclosed in DE A 2125062.\nThe above apparatus for charging a shaft furnace comprises a bin having an inlet and an outlet, the bin being connected through flanges with a top portion of the furnace. The bin outlet is located in a furnace chamber, within its throat zone, and closed by a cone-shaped locking member. The locking member is adapted to move vertically. The furnace throat zone is provided with movable plates which carry distributing plates with a reflecting surface facing the bin outlet. The plates are uniformly circumferentially disposed and adapted to radially move along the guides which are rigidly fixed to the furnace shell.\nIn operation of the apparatus, the locking member is displaced towards the furnace chamber, and a burden is admitted through an annular gap formed to the peripheral region of the throat zone. If the blast-furnace process is disturbed, it is required to change the contour of the charged burden layer. To this end, the plates are radially displaced towards the furnace axis until they reach a position in which the path of the falling burden intersects the surface of the distributing plates, which results in altering the burden flow direction so that the burden moves towards the furnace axis. The above prior art apparatus, however, changes the contour of the burden layer charged only in the peripheral annular region due to the restricted displacement of the movable plate, which reduces the potentials of the apparatus, limiting the burden laying control only to 1/3 the furnace radius.\nMost closely approaching the present invention by the combination of technical features and the achieved result is an apparatus for charging a shaft furnace, comprising a bin for a burden, the bin having an inlet and an outlet, a unit for distributing the burden over a cross section of the furnace, said unit being mounted in a throat zone beneath the bin outlet and adapted to rotate about the furnace axis. The burden distributing unit comprises a horizontal member having at least two guiding members circumferentially disposed around its periphery, each of the guiding members being connected with the horizontal member and consisting of two segments which are arranged sequentially in the direction of the burden flow. A first segment is located in the immediate vicinity to the horizontal member, and a second segment is comprised of at least two plane faces and has a burden shoot edge facing the furnace throat zone (WO Application No. 92/19776, C21B 7/20).\nOwing to the design of the distributing unit, the burden body leaving the bin outlet is divided into at least two flows which form a charge layer in the throat zone with approximately even circumferential grain-size distribution. This is attained due to laying the burden pellets of the same size symmetrically about the furnace axis, thereby significantly reducing the nonuniformity of the circumferential burden weight distribution in the furnace. The uniform circumferential arrangement of the guiding members, with the first segment being disposed in the immediate vicinity to the horizontal member, makes it possible to change the radial direction of the burden flow to the circumferential one. On the guiding member second segments, the circumferential direction of the burden flow changes to the tangential one or towards the furnace axis under the effect of the resultant of frictional, centrifugal and Coriolis forces. This enables the burden to be charged to any throat region and provides the possibility to form a required stockline contour in the throat zone. This is attained by a single type of motion, in particular, by rotation of the burden distributing unit about the furnace axis. The above structural features of the charging unit provide the improvement in the blast-furnace process parameters and the reduction in energy losses. The rotary drive of the unit has a simpler structure.\nIn operation, a burden portion is charged from the burden bin on the horizontal member of the burden distributing unit and further along the guiding members into the furnace chamber. A burden layer is formed in the furnace throat zone by virtue of rotation of the burden distributing unit. Owing to the possibility to vary the rotational speed of said unit, the apparatus provides a desired stockline contour in the shaft furnace. This improves the blast-furnace process characteristics. However, a required stockline contour is not provided when laying the burden by the above apparatus in the region adjacent to the furnace axis, since the non-optimal spatial orientation of the plane faces of the guiding member second segment hinders the admission of a sufficient amount of material into the furnace axial region. As a consequence, the optimum gas distribution over the furnace cross section is not provided, resulting in the increased coke consumption and some reduction in the output. In addition, in the prior art apparatus the position of plane faces of the guiding members fails to be unambiguously defined which makes the manufacture of the apparatus as such problematic, and further hampers the attainment of a desired stockline contour in the furnace throat zone."}
{"text": "This invention relates to apparatus for accurately and uniformly cooling a preformed polymeric sheet. It further relates to apparatus for extruding and cooling a preformed polymeric sheet wherein the gap between adjacent rolls is measured and controlled for accurately and gently cooling the polymeric sheet with minimum stress.\nPolymeric extruded or coextruded sheets are usually extruded out of a slit die of appropriate width, as indicated in the U.S. Pat. Nos. 3,918,865, 3,940,221, 3,959,431, 4,533,510 and 5,466,403, to Frank R. Nissel for example. The hot sheet is then cooled by passing it through a pair or a series of temperature-controlled rolls. Although various roll numbers and arrangements may be used, three rolls are often used, sometimes in planar alignment with each other. The gaps between the rolls are desired to be precisely adjusted according to the desired final sheet thickness. This precision is necessary for a variety of reasons, including elimination of air entrapment between the rolls, which causes adverse or uneven heat transfer or cooling. The rolls typically apply embossing to provide a high quality surface impression on the sheet.\nTraditionally,the sheet die exit is horizontal and the sheet runs through a cooling roll stack at any angle, either upwardly or downwardly. Vertical, horizontal or various angled roll stacks may be used.\nFor various purposes, the art has made efforts from time to time to measure the gap distance between cooling rolls. Gauging the gap distance has been achieved mechanically through the use of feeler gauges. Unfortunately, feeler gauges lack the precision needed for maintaining optimal cooling and surface impression. They provide only a go/no go type reading without giving accurate quantitative measurements. They further tend to scratch the roll surfaces. Feeler gauges also pose the danger of being accidentally drawn into the rolls and permanently damaging them.\nTo avoid these dangers, methods have been introduced where indirect measurement of the gap is taken. This includes measuring the sheet material which passes between the rolls, measuring the distance between the mounting blocks used to support the rolls, and measuring reduced diameter barrel portions. These methods, however, do not and cannot account for actual differences in the diameter of the roll, caused by changes in temperature and other factors. An accurate way to measure the gap between the rolls would be to measure the gap directly, but the reported attempts have not succeeded.\nSeveral devices which monitor a sheet are used to influence the thickness of the sheet product. For example, U.S. Pat. No. 3,809,907 which issued to Schuller et al on May 7, 1974, discloses a device for controlling a bank of viscous material that is formed on, and builds up in front of, the nip of a pair of spaced rolls. Specifically, the rolls are first maintained in close proximity such that the molten material applied to an entry side of the rolls forms a bank of hot material against and before passing through the rolls. The amount of molten steel accumulation is measured by a pair of laser emitters and detectors, located at or near the ends of the rolls. But the laser measures the amount of molten material that is banked up before it flows between the rolls, not the actual roll gap.\nTurning now to the field of plastic sheets, extrusion dies with variable lip portions have been used to form polymeric sheets of varying widths and thickness. Once extruded through an adjustable lip die, the polymeric sheet is substantially at its desired thickness. However, there is a demand for cooling of a preformed polymeric sheet without excessively stressing the polymeric sheet product.\nThere is a need to provide an apparatus for extruding polymeric sheet materials with an extrusion die capable of pre-forming polymeric sheets of accurately predetermined thickness, and a cooling device having cooling rolls defining a gap of precisely known dimension through which the polymeric sheet passes for the purpose of cooling or embossing, wherein the rolls rapidly and uniformly cool the polymer while exerting only minimal pressure on the extruded sheet.\nThis invention includes an extruder and cooling rolls that have an adjustable gap through which a preformed polymeric sheet passes. The gap between the rolls is accurately measured and controlled to allow the rolls to exert only a minimal pressure on the polymeric sheet and to impart the best surface finish to the sheet while subjecting it to the lowest possible stress level."}
{"text": "One type of offshore terminal for mooring a ship includes upper and lower beams extending from the bow of the ship and supporting a turret in rotation about a vertical axis relative to the ship. The lower portion of the turret is anchored to the sea floor, as by heavy chains. An upper bearing arrangement that mounts the turret to the upper beam, includes a thrust bearing that supports the weight of the turret and the load thereon, and that also forms a radial bearing that resists horizontal movement of the turret. A lower bearing arrangement that mounts the turret to the lower beam includes only a radial bearing. The turret is rigid, and the upper and lower bearings must be precisely aligned to avoid large bending stresses on the turret which would result in large radial loads on the bearings that limit their useful life.\nIn practice, it is found difficult to mount the radial bearings on the upper and lower beams to be precisely aligned. For example, even if the bearings are precisely aligned when the turret is installed on the vessel, misalignment will occur at the site where the terminal is installed when heavy chains are attached to the turret. The heavy weight of the chains causes the upper beam to bend, and in bending its outer ends pivot, which results in misalignment of the bearings. A comprehensive analysis can be made to predict the degree of bending of the upper beam, but in practice such analysis is very difficult to perform accurately. Also, it is difficult to install the turret at a shipyard with the required amount of misalignment and consequent bending load on the turret, which is predicted to be compensated for when the chains are installed. A bearing arrangement which avoided the need to accurately predict the amount of bearing misalignment during final installation and the need to install the bearings under high load with high precision, would be of considerable value."}
{"text": "The present invention relates generally to computer system delay elements, and, more particularly, to a system and method for balancing delay of signal communication paths through well voltage adjustment.\nAs technology continues to advance at an unprecedented rate, the transfer of information or data has remained a high priority. In addition, as processing speeds continue to increase, the speed at which information is transferred remains key in increasing overall system speed. Traditionally, parallel busses or cables have been coupled between multiple devices to transfer information during a processing function. However, certain problems that may degrade the transferred data (such as signal loss when transmitted over long distances) have led to improvements in the manner in which such data is transferred. In addition, the multitude of wires or cables necessary for parallel connection may become cumbersome, especially in larger systems. In the past, however, alternatives to parallel connections often led to information bottlenecks, leaving parallel lines as the only viable choice.\nAs the need for bandwidth expands in both existing networks and those now being developed, serial backplanes (or busses) have become an advantageous alternative for solving these problems. Typically, serial backplanes employ a serializer at a transmitting end to convert and transmit data in serial order and a deserializer at a receiving end to convert the data back to parallel form once received. Such high-speed serializer/deserializer (SERDES) receivers have become the benchmark for asynchronous communication and have provided clear advantages over parallel busses. For example, SERDES receivers include transmitters and receivers, and use simplified wiring harnesses (often only a single wire per channel) that typically consume less power than their parallel-coupled counterparts. Higher performance may also be achieved because SERDES receivers reduce the cross talk that often occurs between parallel wires. In addition, SERDES receivers may be employed to transmit data over long distances without the signal degradation experienced with parallel busses, thus ultimately offering increased reliability and fidelity over parallel busses.\nSERDES receivers and transmitters generally classified as mixed-signal devices are comprised of an analog function domain, a digital function domain and a parallel communication path between the analog and digital function domains. The analog domain operates to transmit and receive data over a physical medium or cable, as well as convert between serial and parallel data at high speed, while the digital domain pre or post processes data. Reliable operation of the SERDES device at high speed is dependent not only on the function of the analog and digital domains, but on the parallel communication path that connects these domains. As a result, delay matching of all data path segments that make up the parallel path is required.\nConventionally, the matching of parallel data paths between analog and digital domains has been implemented by using a combination of FET sizing and physical placement techniques in order to reduce the effect of process variance from path to path. However, such methods typically result in the use of larger devices and/or restrictions on the location thereof, both of which hamper physical implementation and can impact circuit speed. Additionally, these techniques do not fully cancel device-to-device variations that can grow, as a percentage of drawn device size, with advancing technologies.\nThe overall performance of an analog/digital design depends, in part, upon the ability to maintain the synchronous tracking of the two domains. Due to broad process variations, and without some means of better adjusting the delays between the two domains, maximum potential performance may not be obtained. Moreover, if the two domains do not track with one another over temperature and time, the tolerance of the overall system will be further reduced. Accordingly, an improved means of real time adjustment of the communication paths between the domains is desirable."}
{"text": "1. Field of the Invention\nThe present invention relates generally to semiconductor manufacturing and related technologies. More particularly, the present invention relates to electron beam lithography.\n2. Description of the Background Art\nAs is well-understood in the art, a lithographic process includes the patterned exposure of a resist so that portions of the resist can be selectively removed to expose underlying areas for selective processing such as by etching, material deposition, implantation and the like. Traditional lithographic processes utilize electromagnetic energy in the form of ultraviolet light for selective exposure of the resist. As an alternative to electromagnetic energy (including x-rays), charged particle beams have been used for high resolution lithographic resist exposure. In particular, electron beams have been used since the low mass of electrons allows relatively accurate control of an electron beam at relatively low power and relatively high speed. The electron beam system can also generate patterns from digital information directly, making the electron beam lithography, maskless lithography. Electron beam lithographic systems include electron-beam direct-write lithography systems and electron-beam projection lithography systems."}
{"text": "Because of the excellent directivity, the power density of laser in the output direction is especially high, which greatly threatens human eyes. The wavelength of lasers used in the existing communication system is generally around 870 nm, 1310 nm and 1550 nm, which are all beyond the wavelength range of visible light. Because it is invisible to the naked eye, the light beam has a greater possibility to hurt human eyes. After the inventions of Erbium-Doped Fiber Amplifier (EDFA) and Wavelength Division Multiplexing (WDM) techniques, the output power of optical signals carried in fibers is even higher than that of conventional Synchronous Digital Hierarchy (SDH) devices. In addition, such apparatuses as RAMAN amplifier may be introduced into the existing communication system, which makes the power of optical signals carried in fibers of some systems reach or even exceed 30 dBm. Such high power of the optical signal greatly threatens the safety of operators and maintenance personnel of telecommunication system.\nThe safety level of laser apparatus and device is defined in Standard IEC60825, which specifies the specific operational measure and safety identifier as well. The International Telecommunication Union-Telecommunication Standardization Sector (ITU-T) constitutes Standard G.664 for laser safety in communication system. According to the Standard G.664, an optical communication apparatus should be able to automatically reduce the output power to a power threshold for safety and even turn off a laser when a laser leakage occurs. The solution is expressed as an Automatic Power Reduction (APR) and an Automatic Power Shutdown (APSD) solution in G.664."}
{"text": "Engine operation can be adjusted to improve engine emissions or fuel economy. However, engine adjustments (e.g., EGR amount and fuel injection timing) that increase fuel economy can also increase engine emissions. One way to improve tailpipe emissions without extremely increasing fuel consumption is to adjust an exhaust gas recirculation (EGR) amount provided to the engine and process engine exhaust with a selective catalyst reduction system (SCR). Increasing the EGR amount supplied to the engine may reduce engine feedgas NOx without extremely increasing fuel consumption. Further, a SCR positioned in an engine exhaust system processing exhaust gas can reduce feedgas NOx to provide lower vehicle tail-pipe NOx emissions. In this way, tailpipe NOx can be decreased via a combination of EGR and exhaust after treatment devices. However, engine fuel economy may decrease somewhat when EGR is increased. As a result, engine fuel economy may be less than is desired.\nThe inventors herein have recognized the above-mentioned disadvantages and have developed a method for operating an engine, comprising: adjusting an EGR amount supplied to an engine in response to an amount of NH3 stored within a SCR catalyst and an amount of urea stored in a tank.\nBy decreasing an EGR amount supplied to an engine when efficiency of a SCR catalyst is high, it may be possible to improve vehicle fuel economy while meeting a desired level of emissions. For example, efficiency of a SCR catalyst may be high when an amount of NH3 stored within the SCR catalyst is greater than a threshold level. As such, the engine may be operated at more fuel efficient conditions via decreasing EGR while tailpipe emissions are maintained below a threshold level by the efficiently operating SCR catalyst. Further, the EGR amount supplied to the engine may be adjusted in response to an amount of urea remaining in a tank when the tank is not substantially empty. Such operation may allow the system to improve engine fuel economy while limiting use of urea to provide a desired interval between urea tank refills.\nThe present description may provide several advantages. Specifically, the approach may reduce engine emissions and improve engine fuel economy. In addition, the approach may also help to provide a desired interval between urea tank refills. Further, the approach may reduce urea injection during periods when SCR catalyst efficiency is high to extend a urea tank refill interval.\nThe above advantages and other advantages, and features of the present description will be readily apparent from the following Detailed Description when taken alone or in connection with the accompanying drawings.\nIt should be understood that the summary above is provided to introduce in simplified form a selection of concepts that are further described in the detailed description. It is not meant to identify key or essential features of the claimed subject matter, the scope of which is defined uniquely by the claims that follow the detailed description. Furthermore, the claimed subject matter is not limited to implementations that solve any disadvantages noted above or in any part of this disclosure."}
{"text": "This invention relates in general to heat exchangers, and in particular to a new and useful heat exchanger for operation at high gas temperatures and high jacket pressures, especially a cracked gas cooler tubular heat exchanger with a tube sheet at the gas inlet side and at the gas outlet side, and with heat exchange tubes arranged in graduated circles.\nTubular heat exchangers of this type include a cylindrical jacket that is sealed off by two tube sheets. The cooling tubes, through which the gas flows parallel to the axis of the jacket are, set into the tube sheets. The flow is around the cooling tubes from the outside. Vertical and horizontal heat exchangers are differentiated according to their positions.\nThe design of the tube sheets presents no particular difficulties. In cracked gas coolers, however, in which the gas inlet temperature is above 800.degree. C. and are subjected to jacket pressures of above 100 bar, the tube sheet on the gas inlet side becomes very hot and therefore has to be made extremely thin to keep the temperature stresses occurring in the sheet small, so that jacket cooling is guaranteed. The protective layer of magnetic forming on the water side during operation, which protects the steel to a great extent against further oxidation, must be retained. Care must therefore be taken that the deformations are kept within limits in spite of high stress to which the tube sheet is exposed, so that this protective magnetic layer is not destroyed.\nIn cracked gas coolers that have to operate under the aforementioned conditions, and which have been designed with very thin \"hot\" tube sheets because of the temperature stresses that occur, in order to be able to absorb the loads, tie bolts have been welded to the so-called hot tube sheet perpendicular to its central plane. The bolts are distributed uniformly over the tube sheet surface in tube gases passing through, which were fastened at the other end to very rigid cross bars that are supported on appropriately heavy duty forged rings on the tank jacket.\nThe so-called cold tube sheet of the cracked gas cooler had thick walls corresponding to its loads. Design calculation on this basis is very tedious."}
{"text": "The present invention relates to an injection molding machine having an injection device and a mold opening and closing device, particularly to a motor-driven injection molding machine using motors as drive sources in the injection device and the mold opening and closing device.\nAs is well known, according to a motor-driven injection molding machine, respective functions of injection, mold opening and closing, measurement and ejector, are realized by individually installed motors. Induction motors are normally used for such motors. An explanation will be given of an example of a motor drive for such a motor in reference to FIG. 1. In FIG. 1, a motor drive 10 is for a three phase induction motor IM and includes an inverter section 20 and a rectifying circuit section 30. The inverter section 20 includes a plurality of switching elements 21 for carrying out switching operation for respective phases of three phases. The rectifying circuit section 30 is realized by a condenser input type diode rectifying circuit. The diode rectifying circuit includes a plurality of diodes 31 for carrying out rectifying operation for respective phases of a three phase power source PS.\nMeanwhile, the induction motor IM for the injection device or the mold opening and closing device generates regenerated electric power in deceleration thereof. Heretofore, in the motor drive, regenerated electric power of the induction motor IM is dissipated as heat by a resistor 41 provided to a dynamic brake circuit section 40. That is, when the induction motor IM generates regenerated electric power, the switching element 42 is switched on and the regenerated electric power is dissipated by the resistor 41. From a view point of reduction of electric energy, wasteful electric energy is dissipated. In addition thereto, there is constituted a heat radiating device for radiating generated heat of the resistor 41 and therefore, there poses a problem of an increase in dimensions.\nAs a measure with regard to the above-described regenerated electric power, there is provided even a system enabling power source regeneration by replacing the rectifying circuit section 30 by a power source regenerating converter. However, the power source regenerating converter is expensive.\nFurther, other than the above-described, it is the current situation that heretofore, according to the injection molding machine, no measure has been taken for electricity interruption except a necessary minimum power source backup function. Therefore, when molding operation is stopped by electricity interruption in the midst of a molding cycle during the operation, in the case in which resin is not sufficiently charged in a mold, an ejector mechanism cannot eject unmolded product in the mold. In this case, the unmolded product, that is, a failed molded product must be taken out by manual operation. Time and labor is required in taking it out and depending on cases, the mold must be disassembled. Further, when the operation is stopped in the midst of the molding cycle by electricity interruption, much time and labor is required and skill is needed in resetting peripheral apparatus of the injection molding machine in restarting thereof.\nHence, it is an object of the present invention to provide an injection molding machine capable of reducing power consumption by achieving effective utilization of regenerated electric power of a motor.\nIt is another object of the present invention to provide an injection molding machine capable of utilizing regenerated electric power of a motor to operation in electricity interruption.\nAn injection molding machine according to the present invention comprises a drive source comprising at least one motor and a motor drive, a charger for storing regenerated electric power of the motor and electric power from a power source, a charging and discharging circuit connected to the charger for charging electric power to the charger and for discharging electric power therefrom, and a controller for controlling the charger to charge electric power and for controlling the charger to supply electric power to the motor. The controller controls the charger by controlling the charging and discharging circuit based upon an operating condition of the injection molding machine.\nAccording to another aspect of the present invention, a method of providing drive power to an injection molding machine is provided. The method comprises the steps of controlling a charger to charge electric power from a power source and to supply electric power to a motor of the injection molding machine. A controller provides control signals to a charging and discharging circuit based upon an operating condition of the injection molding machine."}
{"text": "Robot welders are automated welding machines which operate under the control of a microprocessor. They are being used in more and more applications because they provide several advantages with respect to human welders. Some, but not all, of the advantages are that they do not need to take breaks, they can tolerate adverse environmental conditions, they can work in very tight areas, and they yield a consistent, quality weld.\nHowever, a robot, like a human, must be told what to weld. A program for a robot welder can be written which specifies the exact spot or line that must be welded. Typically, an operator specifies the starting point of a weld and the length of the weld or the ending point of the weld. This information is given in precise coordinates, for example, start at point x=1.05 inches, y=1.77 inches, and z=0.0 inches and end at x=2.00 inches, y=2.00 inches, and z=0.5 inches. In this case the robot will weld on a seam defined by a line between the given starting and ending points. This procedure is simple and straightforward, but is very time consuming because the operator must determine the starting and ending points of the weld. This procedure is also prone to errors because the operator must manually enter these starting and ending points into the program and the operator may make errors in measuring the points or in entering the points into the program. Further, some points may be difficult or impossible for the operator to measure manually.\nThe robots have sensing devices, often referred to as touch sensors, which indicate that contact has been made with an object. In welding environments, the workpiece is typically a conductive material, such as iron or steel. Therefore, in these environments, the touch sensor is often the tip of the welding torch, or the welding wire, and the robot senses when the tip or wire has reached the workpiece by simply monitoring the output voltage or current at the tip or wire. When contact is made, the output voltage will drop and the output current will increase. The operator can therefore operate the robot to determine the coordinates of the starting and ending points of the weld, and then enter these coordinates into a welding program. This procedure is simple and straightforward, but is also very time consuming and prone to errors because, even though the robot makes the measurements precisely, the operator may make errors in entering the measurements.\nThus, each weld requires a substantial amount of attention by the operator. In some cases, the amount of time required to determine the coordinates of the weld and enter the points into the welding program can easily exceed the amount of time required for the robot to perform the weld. Also, because of human measurement and data entry, errors are quite possible. Therefore, there is a need for a method of operating a robot welder so as to reduce or eliminate the amount of attention required from the operator.\nFurther, once a welding program is written it will generally be useful for only one specified weld. That is, as long as the coordinates of the welds are identical (same starting and ending points), then the robot will faithfully repeat the weld. This is most useful on assembly lines where the robot performs the same welding operation on each workpiece which comes down the assembly line. However, if the coordinates of the welds are not identical then a different program is required. This may occur where one robot is required to make several welds on the same workpiece. Thus, the amount of memory required to store the several programs necessary for multiple welding operations can be substantial or even beyond the capability of the robot.\nRobot welders, although useful, do only what the programmer instructs them to do. This is of little concern if there is only one robot performing a welding operation because, other than hitting the workpiece, there is little danger of the robot hitting anything. However, speed of assembly is frequently a customer requirement so it would be beneficial if two or more robots could be assigned to weld on the same workpiece at the same time. The problem with two or more robot welders is that they may try to occupy the same space at the same time. Two humans assigned to perform conflicting welding operations may be able to discuss and resolve the problem among themselves. They may also contact a supervisor to determine whose work will be done first. However, robots do not have the capability to recognize that a problem exists, much less discuss and resolve the problem or seek the guidance of a superior authority. Therefore, each robot follows its own program, mindless of the presence of any other robot. The result at some point, with some programs, will be a collision. The resulting damage to the robots and possibly even the workpiece can be extensive and expensive, and can cause serious delays.\nIt is difficult, if not impossible, to know in advance of the performance of an operation whether or not there will be a collision between the two robots. When programming a robot to perform a welding operation, the operator may specify a series of points to which the robot must move, with or without specifying exactly how the robot will get there. However, the microprocessor which controls the robot has its own set of manufacturer-installed operating instructions as to how the robot should operate. This problem is compounded by the fact that a robot typically has numerous degrees of freedom. For example, the Panasonic industrial robot model AW-010A has 6 degrees of freedom (rotation, upper arm, front arm, rotating wrist, bending wrist, twisting wrist), and two directions of linear motion (forward-backward, up-down). The particular moves generated by these operating instructions are frequently not the same moves, or in the same sequence, that the operator would have thought to use to accomplish the same result. Therefore, even if the operator specifies that a robot is to go from one point to another point, the operator has no information on what motions the robot may implement to do so. Further, the operator may have to particularly specify parts of the sequence to get the welding torch in the right position and at the right angle while avoiding parts of the workpiece, for example, an overhanging flange from a T-beam.\nThis problem is made worse if one attempts to simultaneously operate more than one robot on the same workpiece. None of the robots knows, or has the capability to know, where the other robots are located, what the other robots are doing, what the other robots are preparing to do, or how the other robots are going to accomplish whatever their next actions may be. Therefore, in the past, the only certain ways to avoid a collision were to put only a single robot on a workpiece, or to space the robots far enough apart to where it was impossible for one robot to intrude into the operating area of another robot. However, a single robot provides for slow operation. Further, in some situations, the small size of the workpiece may prevent spacing the robots such that they cannot hit each other. Also, operating space is like any other resource; it has a value and should be conserved whenever possible. There is therefore a need for a method which provides for simultaneous operation of the robots in close proximity to each other while preventing the collisions which, with the current state of the art, are probable.\nA sophisticated spatial analysis computer program could be used to generate a computer model of each robot and the workpiece and perform a mathematical computation to determine if any robot will attempt to occupy the same space as any other robot at the same time. However, such a program is expensive and requires a substantial amount of time for programming of the dimensions and variables. Such a program is also very computation-intensive and requires a large amount of memory. This may tie up the resources of a company which has its own computer, or increase the computer time rental costs for a company. Furthermore, a spatial analysis program may not have been written for the particular processor being used to control a robot. In addition, there may be inadequate time, at a processing speed which is not cost prohibitive, to perform the calculations in real time. Also, any change to any program may require running the analysis again to verify that the change has not caused a collision.\nWelding robots typically have current sensors associated with each driving motor. Therefore, if the robot contacts an object and a driving motor stalls out, then the current will increase. This feature is often referred to as an amp-out condition and can be used to determine if the robot has hit an object or another robot. However, at the speeds at which robots operate, the damage will still be serious. It is, of course, possible to program the robots to move at a very slow rate so that the amp-out condition can be used to determine if a collision is occurring and to stop movement of the robots so that the collision will not cause further damage. However, this only prevents further damage, it does not prevent any damage caused by the initial contact nor resolve the original problem: both robots trying to be in the same space at the same time. Also, by moving at this slow speed, the time needed for a welding operation is greatly increased. Further, the slow movement approach is only useful for positioning. A welding operation must usually occur at a predetermined speed. Moving the robot arm at a slower speed during a welding operation may result in the workpiece being damaged by excessive heat. The workpiece may be warped or even burned through, or an excessively large welding bead may be deposited. Therefore, there is a need for a method of operating robot welders in a manner which avoids collisions between the robots.\nBefore beginning a welding operation, the robot must know how to get the torch into the compartment in which the welding operation is to be performed. This can be done by the operator, by manually measuring the coordinates, or by manually positioning the robot in the compartment. However, manual measurement of the coordinates is time prone to errors, and manual positioning is time consuming. In many cases, the workpiece is designed using a computer aided design (CAD) program. This CAD program contains information which could be used to determine the entry point into a compartment for the torch. This would eliminate the need for manual measurement or control, would improve speed, and would eliminate errors. Therefore, there is a need for a method for specifying a compartment entry point for a welding operation based upon data provided from a CAD program.\nEven after the torch is guided into the compartment in which a welding operation is to be performed, the starting and ending points of the weld must be determined. This can be done by the operator, by manually measuring the coordinates, or by manually positioning the torch to determine the coordinates, and then loading these coordinates into the robot. However, these procedures are time consuming and may introduce errors. Some robots have a \"touchsensing\" capability. That is, the ability to locate the boundaries or walls which limit the operation area of the robot. This capability could be used to define the coordinates of the weld. Therefore, there is a need for a method for automatically determining the starting and ending coordinates of a weld. There is also a need for a method for automatically determining the starting and ending coordinates of a weld using the touch-sensing capability of the robot.\nProgramming a robot by going out on the shop floor and using a pendant to teach the robot the steps necessary for the welding operation is useful, but is time consuming and, therefore, inefficient and expensive. Further, it requires the operator to go into an environment which is not climate controlled. If the weather is particularly adverse then the time that the operator can spend on the shop floor may be very limited. To circumvent this problem, some facilities use programmers to write the welding program by looking at the CAD design, rather than by going out on the shop floor. The programmer is thus in a climate-controlled environment and is not affected by the weather. However, this approach is often unsuccessful on the first attempt at programming. This occurs because there are limits on the agility of a robot. The robot may be able to turn its wrist through most of, all of, or even slightly more than, a 360 degree arc but there will be a limit. This limit condition is installed in the robot by the manufacturer and prevents excessive turning so that wires, hoses, and cables do not get wrapped around the robot arm and/or broken. If this limit is implemented by software it is frequently referred to as a soft limit condition and if it is implemented by hardware, such as a switch, it is frequently referred to as a hard limit condition. If a program instructs the robot to move more than the soft limit condition allows, the soft/hard limit causes the robot to disregard the program instruction. A programmer may unknowingly write a program which attempts to cause the robot to exceed a limit condition. When the offending program instruction is reached then the robot simply stops operating. The programmer may have to run the program several times to determine why the robot stopped. Then, the programmer will have to modify or completely rewrite the program in order to avoid reaching the limit switch condition. In some cases, the programmer may have to modify or rewrite the program by going out on the shop floor and using the pendant. There is therefore a need for a method of generating a welding program which avoids the limit switch problem and which also prevents the operator from having to go to the shop floor for the programming.\nBefore the welding operation is performed, the angle of attack of the torch of the robot must be specified. There is an optimum angle of attack for most welding operations. However, merely specifying that angle in the welding job for a robot is not useful because the compartment in which the welding is to be done may be too deep, or have a beam with too much of a flange overhang, to obtain the desired angle of attack and the robot may merely jam the torch into the workpiece in a vain attempt to achieve the desired angle of attack. Therefore, there is a need for a method for automatically determining the positioning of a torch to provide for an optimum angle of attack for a welding operation.\nFurther, the optimum angle of attack for the torch is not fixed, even for a known compartment design. If the orientation (pitch or roll) of the compartment is changed then the optimum angle of attack for the torch will vary to compensate for the tendency of the molten metal to flow downhill. Also, if one takes the program for a compartment which has no pitch (inclination) and tries to use the program for a compartment which has some pitch, the robot may be able to get into position to start the weld but may encounter a limit condition as the robot attempts to perform the weld. Thus, a different orientation (yaw, pitch, roll) may require the robot to use a different approach path in order to get the torch into the compartment and be able to complete the weld. Therefore, there is a need for a method for generating a program which is responsive to the orientation of the workpiece.\nOnce the robot has successfully placed the torch within the desired compartment and the torch is positioned at the best available angle of attack, the robot must move the torch along the joint of the pieces to be welded together. Point-to-point specifications are very useful if the joint is along a straight line. However, if the joint is along a curve or bend, or there is a change in the direction of the joint, then point-to-point specifications become less useful. Attempting to create a curve by a series of point-to-point specifications is time consuming and also prone to errors because the ending point of one path must precisely correspond to the starting point of the next path. Also, if the curve is severe, then numerous point-to-point specifications may be necessary to simulate the curve closely enough to keep the torch on the joint. There exists a method for automatically tracking the joint for a welding seam where the seam is a joint between two planar surfaces. However, in some cases, the weld to be made is not a joint between two planar surfaces but is a joint between two perpendicular surfaces. The prior art seam tracking method frequently fails to provide the desired results in the case of perpendicular surfaces. Therefore, there is a need for a method for automatically tracking the seam between two perpendicular surfaces.\nIn some cases, the weld to be made is not a joint between two surfaces but, rather, is simply a welding bead placed along the edge of a piece to eliminate the rough edge, or provide for rust prevention, or provide a better edge for a later welding operation, etc. As in the case of a seam, point-to-point specifications are very useful if the edge is a straight line. However, if the edge is curved or bent, or there is a change in the direction of the edge, then point-to-point specifications become less useful. Attempting to create a curve by a series of point-to-point specifications is time consuming and also prone to errors because the ending point of one path must precisely correspond to the starting point of the next path. Also, if the curve is severe, then numerous point-to-point specifications may be necessary to simulate the curve closely enough to keep the torch on the edge. Therefore, there is a need for a method for automatically tracking the edge of a component so that a welding bead can be applied to the edge."}
{"text": "1. Technical Field\nThe present invention relates to a webbing winding device that constitutes a safety belt device of a vehicle.\n2. Related Art\nIn a so-called vehicle seat belt retractor disclosed in Japanese Utility Model Application Laid-Open (JU-A) No. H06-8122, when a winding shaft rapidly rotates in an extracting direction due to a tension applied to a webbing, a latch cup also rotates in the extracting direction because a lock portion of a lock member meshes with an internal tooth of a latch cup. As such, when the latch cup rotates in such way, an inner peripheral portion of an elongated hole-shaped pole guide formed at the latch cup presses a pole guide protrusion so as to rotate a pole. The rotating pole allows the engagement portion thereof to engage with a latch plate provided in the winding shaft, thereby regulating the rotation of the winding shaft in the extracting direction.\nIn the above-described configuration in which a protrusion such as the pole guide protrusion is inserted into an elongated hole such as the pole guide hole and an inner peripheral portion of the elongated hole presses an outer peripheral portion of the protrusion portion so as to transmit the movement, it sometimes happens that, for example, if either one provided with the elongated hole (the latch cup of JU-A No. H06-8122) is inclined in a manner such that the opening direction of the elongated hole is inclined with respect to the protruding direction of the protrusion portion, the inner peripheral portion of the elongated hole presses the front end of the protrusion and thus a force acts in a direction in which the one provided with the protrusion (the pole of JU-A No. H06-8122) is forcedly inclined.\nIn this way, there are cases that the movement may be transmitted from one provided with the elongated hole (the latch cup of JU-A No. H06-8122) to the other one provided with the protrusion (the pole of JU-A No. H06-8122) in the state in which the one provided with the elongated hole and the other one provided with the protrusion are inclined."}
{"text": "1. Field of the Invention\nThe present invention relates generally to chimneys. More specifically, the invention is a funnel attachment for a chimney sweep vacuum which removes soot from a fireplace flue or work space without the need for tarps and the like to catch excavated soot spillage.\n2. Description of the Related Art\nThe general use of particle removing apparatuses have been in existence as early as 1920. It was A. E. Burges (U.S. Pat. No. 1,371,468) who was credited with the idea of harvesting or removing cotton in its natural state by a suction based apparatus. This principle took root in the minds of others such as E. Gray (U.S. Pat. No. 1,465,711) and Engstrom (U.S. Pat. No. 1,979,873) which utilized suction channels for removing soot from furnaces in 1922 and 1933, respectively. As one might expect, these systems where often quite cumbersome to operate, and often include fixed mechanical attachments for limited use to clean up residual debris in difficult to reach areas. U.S. Pat. No. 3,795,181 issued to Lawson discloses a flue cleaning device as a fixed attachment to a flue housing. As described therein, a turbine is used having a plurality of discharge jets which discharge fluid within the housing to remove particles therefrom.\nLater developments by Evans (U.S. Pat. No. 4,807.590) and Broussard (US) sought to provide a vacuum system as a portable device, but were met with mechanical limitations which inhibit the effective use of vacuum heads and the like in and around flue mechanisms of various sorts. Thus, a chimney vacuum system which is not messy, easy to use and which can be used in and around various flue mechanism or fixtures uninhibited is needed. Other patents granted to Clarkson et al. (GB 421,664), Newell (GB 622,943), Neuman (CA 1 181 55), Holter (DE 3529310), Seiji (JP 3102112), Asaph (CA 2010630) and Kupracz (CA 2040616) provided conventional chimney vacuum features which were face with the same limitations recited above, and thus are considered to be generally relevant to the invention as herein described.\nNone of the above inventions and patents, taken either singularly or in combination, is seen to describe the instant invention as claimed. Thus a chimney vacuum solving the aforementioned problems is desired.\nThe chimney vacuum system according to the invention includes an adapter head and hose assembly for quick and mess free removal of soot from interior rooms in residential and commercial buildings. The adapter is substantially conically shaped and has a neck portion for attaching an effluent or extension hose thereto. The shape of the adapter head is defined by an outer diameter which tapers in volume to an inner diameter which terminates a predetermined length to define a cylindrical channel or neck. The area defined between the outer diameter end and the inner diameter end forms the body of the adapter. Defined therein is a sealing gasket which forms conic section of predetermined area and distinct material properties (i.e rigidity, flexibility and expandability) from the body portion of the adapter head for providing a mechanical and fluid seal between at least one flue mechanism element such as a flue handle or the like.\nThe adapter is mechanically secured to an effluent hose at the neck for transporting contained flue soot under a vacuum from a work space to a remote location. The adapter is made as a single unitary structure with transparent plastic features for visually identifying soot removal. An excavation brush is used as another elemental feature of the system having tines or brush fibers which span the entire area of the flue to excavate and flush accumulated soot from the interior walls of the flue at selective heights within the chimney.\nAccordingly, it is a principal object of the invention to provide a chimney vacuum system for removing accumulated soot from flues.\nIt is another object of the invention to provide a chimney vacuum system having an adapter which accommodates at least one chimney or flue mechanism without inhibiting an effective mechanical and fluid seal for removing accumulated soot from flues.\nIt is a further object of the invention to provide a chimney vacuum system which is simple to operate and deploy within interior rooms residential and commercial buildings.\nStill another object of the invention is to provide a chimney vacuum system which minimizes dust and soot residue in a work space.\nIt is an object of the invention to provide improved elements and arrangements thereof for the purposes described which is inexpensive, dependable and fully effective in accomplishing its intended purposes."}
{"text": "1. Field of Invention\nThe subject invention is generally related to safety and surveillance equipment for aircraft, ships and other commercial vehicles and is specifically directed to a comprehensive multi-media security surveillance system for same while in port or terminal and/or unattended whether taxiing or parked or docked. The system of the subject invention also permits tracking while en route.\n2. Discussion of the Prior Art\nSecurity is of ever increasing importance. This is particularly true with respect to commercial airlines as more and more people and freight are moved in this manner and aircraft and other commercial transports increasingly become the targets of terrorists activities. The airways are becoming increasingly crowded with traffic. Global tracking systems are now in place to monitor the flight of the aircraft from the moment it lifts off until it safely lands at its destination. Radar and navigational positioning systems are commonplace both on the aircraft and at the ground tracking stations. All of these electronic systems have increased the overall safety record of commercial traffic to new standards as the number of miles flown continues to escalate.\nIn addition, the on board avionics including electronic monitoring and diagnostic equipment, particularly on large commercial jets, continues to evolve, giving both the on board crew and the ground assets more complete, accurate and up to date information regarding the condition of the aircraft while in flight. Flight recorders long have been incorporated in order to provide a record of each flight and in order to provide critical information to aid in the determination of the causes of an accident or malfunction should one occur.\nHowever, one area which has been neglected with the ever increasing availability of electronic surveillance is the security of the aircraft or other vehicles or vessels, including, but not limited to, over-the-road vehicles, ships and other commercial vehicles (collectively referred to as commercial transports), particularly when unattended. Typically, when an aircraft is on the ground, or in port, and unattended the only security provided is the security of the location. If the security of the area in which the commercial transport is stored is breached, the commercial transport is an easy target. In most cases, even the access doors are left open and further, for obvious safety reasons, are designed not to be locked from the outside. Many critical areas of the commercial transport are left exposed such as in an aircraft, by way of example, the landing gear, the engine housing and critical wing and tail components.\nWith terrorism and sabotage an increasing problem there is significant need to develop an integrated system capable of providing good physical/visual and/or audio surveillance as well as monitoring of the environmental, security and motion conditions of the commercial transport and various components while the commercial transport is on the ground. For example, a good visual surveillance system would give instant evidence of a breach of commercial transport security, could sound an alarm and could immediately secure the area. Such a system would also permit the recording of visual information to provide a history for later review, providing yet another source of information for increasing the overall security of commercial transportation.\nWhile such a system would be of great benefit to the commercial transport and airline industries in general and to the commercial airlines in particular, there are no integrated systems currently available which adequately meets these needs.\nThe subject invention is directed to a comprehensive multi-media safety, tracking and/or surveillance system, which in the preferred form provides both visual and/or audio information as well as critical data such as location, direction, intrusion, fire and/or smoke detection and/or status of environmental conditions and/or transport systems status. In my aforementioned patent and copending applications, incorporated herein by reference, detection and sensor systems are utilized to provide the flight crew and/or a ground tracking station for commercial aircraft critical information during flight and/or to record the information and data generated during flight for later reconstruction of catastrophic events. The subject invention is an expansion of this concept and adds not only ground security and surveillance, but tracking while in port or on the ground as well as while in route, as well as incorporating the onboard systems of the aforementioned patent and applications. It is an important feature of the invention that the transmitting network for provides a comprehensive communications link among the transport, per se, stationary and mobile stations on the ground, as well as personal security units, with the craft or vehicle being monitored by both strategic sensors onboard the commercial transport and also strategically placed sensors on the ground.\nIn the preferred embodiment of the invention, a wireless LAN (local area network), WAN (wide area network) and/or other wireless transmission scheme is used as the transmission system of choice and permits the use of IP protocol for providing a comprehensive, multi-media security and surveillance system with communication and command capability between the transport and the ground station, the transport and personnel, personnel and ground and personnel to personnel, including both mobile and permanent ground based stations. Monitoring of and communication with support personnel and vehicles as well as between transports is also provided. A digital wireless voice intercom is provided for security purposes and for communication between the onboard crew and the ground based personnel. In the preferred embodiment video intercom is also provided. Digital wireless (such as, by way of example, LAN) based file communication capability permits the transmission of information such as route or flight plans or gate and dock information. As an example, a LAN or WAN has worldwide tracking capability adapted to be used in connection with a global satellite communication system such as IRIDIUM, wherein the entire path and status of the commercial transport may be monitored while airborne over satellite connections.\nOne important feature of the invention is the ability to remotely monitor a commercial transport while on the ground, whether or not the commercial transport is attended. This will permit detection of unexpected events, breach of security, change in environmental conditions and other activities both on and in the vicinity of the commercial transport. A G.P.S. system may be included to provide accurate positioning information of the commercial transport, establishing the parked position of the commercial transport at any time, as well as tracking its movements.\nIn its preferred form, a plurality of sensor units, which may include at least one video or image sensor/device and/or at least one audio sensor and/or at least one motion sensor, are placed strategically about the interior and exterior of the commercial transport and at strategic ground based locations. In addition, strategically placed motion detectors, fire sensors, smoke sensors, door or latch sensors and other monitoring equipment are incorporated in the system. A comprehensive system incorporating these various sensing devices provide a broad based, multi-media safety, security and surveillance system for monitoring commercial transports at any time, whether or not attended.\nIn addition to safety and/or surveillance issues, the comprehensive data collection scheme of the subject invention provides a system permitting enhanced monitoring and/or response to crew generated work orders or re-supply orders, and may even avoid the requirement that the crew order certain supplies. For example, by monitoring the fuel, fresh water, waste water and/or hydraulic levels onboard and transmitting this to a ground station, refueling, water delivery and/or hydraulic fluid check and supply may be initiated by the station crew and prepared for delivery when the commercial transport arrives in port. The performance parameters of the commercial transport may also be monitored and may be utilized for initiating maintenance procedures, for example, even before the commercial transport is in port. Pre-flight or pre-mission checklists may be enhanced or automated by monitoring the critical functions and criteria via the system of the subject invention. The system of the subject invention greatly enhances maintenance procedures and efficiency. Where desired, the system is capable of permitting the commercial transport to carry its detailed maintenance record onboard, permitting full access to such information at remote locations. The maintenance record can be routinely updated or polled from the home based maintenance station using the system\"\"s unique uplink capability. The ability to both send and receive information will support remote control of the commercial transport onboard systems such as lighting, strobes, alarm setting/resetting, environmental controls, locking systems, siren or other audible signals, fuel flow, fire detection and the like.\nOnce in port, the system of the subject invention permits complete monitoring of on ground movement, and allows the monitoring of other commercial transport in the area to assure that the various transports do not interfere with one another. This provides collision avoidance, and can be utilized both on the ground and in the air or in route via water or land. Current airborne collision avoidance is accomplished by use of a radar transponder. Aircraft position is located by radar xe2x80x9cechoxe2x80x9d response and altitude by a xe2x80x9creporting altimeterxe2x80x9d reading being returned to the radar system encoded in the transporter return. Use of a satellite based LAN or WAN will provide an xe2x80x9cintranet in the skyxe2x80x9d, providing much more accurate G.P.S. position, altitude, heading, speed and other navigational information to the FAA and other operators and computer tracking and monitoring stations, thus enhancing collision avoidance information.\nSituational awareness is also provided by the subject invention. All transports in the terminal area are provided with a G.P.S. location sensor such that the home or ground crew will be able to track and identify the location of every transport in the terminal. This provides better flow of the commercial transports in the terminal area, assuring that proper distance is maintained and appropriate pathways are followed.\nThe system also permits full situational awareness capability where all ground or water transports in the are provided with G.P.S. location information such that the ground crew will know where all assets are at any point in time. This can provide both collision avoidance as described and also check to assure that the transports are in an authorized area. A composite of all transport location information can be used to provide a xe2x80x9clivexe2x80x9d display of all assets in the area. Logging of this information will provide good archival information in the event a reconstruction of events, such as a security breach or collision, is required.\nThe comprehensive multi-media system of the subject invention permits the collection and dissemination of virtually all data associated with the commercial transport at any time, both while in port or in service. In the preferred embodiment a combination of sensors systems are used, with sensors being installed within the transport, on its exterior and at ground-based locations for monitoring the transport when is in port. In such areas where ground based systems are not available, the transport-installed systems still provide useful and enhanced information over the prior art. Likewise, in those areas where unequipped transports enter a system equipped port, the ground based system of the subject invention can communicate via standard ground-to-air radio to provide useful information such as perimeter surveillance and the like. For example, even without the use of on-board systems, the identification number (such as the tail number on an aircraft), owner, state or country of origin and other identifying information can be matched with available data to provide immediate and accurate identification of a specific commercial transport. This permits efficient tracking and response capability of the transport in port, on the ground, or anywhere in the world using satellite communications.\nIn the preferred embodiment, fixed view and steerable video cameras may be incorporated either on the commercial transport or independently of the transport at ground based sites where commercial transport is located in order to monitor movements around the perimeter of the monitored commercial transport. It is also desirable to include focusing and/or timing functions so that selective pan, tilt and/or zoom (x,y,z) positioning can be utilized. The cameras may be activated and/or aimed and/or focused based on the location data provided by a G.P.S. system integral to the monitored commercial transport, may automatically pan an area, or may be manually operated by crew or ground personnel. Automatic tracking of each transport in the terminal by one or more tracking cameras in conjunction with a recording device can provide an archival record of each asset in case of a detrimental event, such as fire, terrorist event, theft, collision and the like. Information can also be supplied to assets that are not G.P.S. equipped.\nSeveral video cameras may be placed such that the lens of each is aimed through a window opening provided in the fuselage or body in order to provide video imaging of the engines, tail section, and/or landing gear and other functional components of an aircraft. Cameras may be placed throughout the interior of the commercial transport on the flight deck, in the cargo hold, in passenger cabin and/or other desired spaces including on the ground outside the commercial transport. The audio sensors/transducers and/or other sensors and detectors are also strategically located throughout the commercial transport and positioned at strategic locations both internal and external of the fuselage. External sensors based on the ground area surrounding the commercial transport may also be added.\nIn its simplest form, current sensors are already on the commercial transport coupled with strategically based ground sensors and may be used to provide surveillance and/or warning system. Thus, a basic system may be implemented with a minimum of alteration to the commercial transport and a minimum of expense.\nWithin the commercial transport, the system may be hardwired or may use wireless transmission and receiving systems. The wireless system is particularly useful for adapting the system as a retrofit on existing equipment and also provides assurances against disruption of data transmission during structural catastrophes such as fire or airframe breakup. In the preferred embodiment, the wireless system is fully self-contained with each sensor unit having an independent power supply and where appropriate, a sensor light source. The ground sensors may likewise be hardwired or use wireless transmission and receiving of video and/or alarm telemetry signals. The ground security system may include motion sensitive, weight sensitive, infrared sensitive, audio sensitive, or other typed activation system so that the equipment is not activated until some event is detected, i.e., the system is action triggered. The ground communications link, monitoring and/or recording systems for collecting and/or transmitting the data as disclosed in my cozening applications may be adapted for processing the information gathered by the on-ground security system and, in the preferred embodiment. The wireless system may use radio frequency transmission and may incorporate the wireless communication system already in place as an integral component of the system. Where desired, a wireless local area network (LAN) or other wireless system may also be utilized for intercommunication among the system components. Preferably, the entire capture, retrieval, monitor and archive system is installed utilizing the wireless transmitting/receiving system in order to assure that transmission will not be lost in the event of a power shutdown or a failure causing possible open or shorted circuit conditions which could occur in a hard wired system.\nA commercial transport equipped with the ground surveillance system of the subject invention may not always be located at a port or terminal equipped with a ground security system. In the preferred embodiment of the invention, the on-board system is self-contained and can operate on a stand-alone basis at sites where compatible comprehensive electronic ground security is not available. In those sites with a compatible ground surveillance system, the on-board system communicates with the site-based system to provide information to airport ground personnel and security personnel. The system of the present invention also lends well to a deployable surveillance device carried by the transport, which can be deployed at unequipped sites to permit off-craft monitoring while the commercial transport is at the port or terminal. The system can be positioned at a strategic location within the site whenever the commercial transport is unattended to permit monitoring of the commercial transport from a remote location. The deployable device is then retrieved and stowed in the commercial transport when the commercial transport departs from the site.\nIn the preferred embodiment, the system will transmit any detected information to a monitor system located at a ground control security station, typically located somewhere within the terminal, tower and/or safety sites such as security stations and fire stations. Detection of activity or fire can sound local and/or remote alarms and/or dial emergency numbers. The data may also be recorded on the standard recorders provided onboard the commercial transport and/or on ground based recorders of conventional type, digital type or a computer based logging system. The security station has instant live access to all of the image and/or audio signals as they are captured by the sensors, and where used, the commercial transport recorder will make an historic record of the images for archive purposes. Where random access recording techniques are used, such as, by way of example, digital random access memory storage devices, the information by be readily searched for stored information.\nIf unauthorized personnel breaches the security area and the audio and/or video equipment is activated, signals will be immediately transmitted to the security station. This will give immediate access to information identifying the activity and the personnel involved. Further, in the preferred embodiment of the invention, an alarm system will be activated for securing the immediate area and taking counter measures to tighten security such as remote operation of lights and doors, and respond to a breach of same.\nIn the one embodiment, information from the plurality of sensors on the transport is synchronized through an on board capture/multiplexing system whereby the plurality of data, including visual image data, may be displayed, recorded, and/or transmitted in either a split screen or serial fashion. A xe2x80x9ctime-stampxe2x80x9d or chronology signal may also be incorporated in the data scheme. Any signal which is capable of being captured and stored may be monitored in this manner. Utilizing the wireless system of the invention in combination with the battery back-up power supply, it is possible to continue collecting information without using ground power or commercial transport power. This assures that the system will operate even if power is disrupted for any reason such as, by way of example, tampering by unauthorized personnel or by fire. In its simplest form, only triggered (activated) sensors are active, i.e., an activity at the site causes a triggering effect and activates the sensor, and only the signals generated thereby are transmitted to the security station. In such a system, multiplexing of continuous signals is not nearly as critical. The xe2x80x9ctime-stampxe2x80x9d is particularly useful as an aid in reconstructing the events in a xe2x80x9cpost-eventxe2x80x9d investigation.\nIn the one embodiment, the system includes a plurality of strategically located video image sensors and/or audio sensors, each sensor adapted for transmitting the signals to a multiplexer for distributing the signals to monitors and/or archival recorders. The data multiplexer combines all of the signals from the various detector circuits to provide a data stream suitable for transmission over the wireless system.\nThe LAN transceiver is the interface into the LAN. The LAN transceiver can accept software downloads from various system elements to enable the multi-media sensor system to be maintained or upgraded to perform other functions. Other sensors may also be incorporated in the system, such as motion sensors, smoke and/or fire sensors and the like. The system is configured for selectively transmitting all of the data on a xe2x80x9creal-timexe2x80x9d or xe2x80x9cnear real-timexe2x80x9d basis, i.e., the data is delivered with only delays for processing time such as compression/decompression, multiplexing and the like. The system is also adapted to provide the monitors access to serial, synchronized full screen view of each of the cameras, in sequential viewing, or alternatively to provide split screen or multi-monitor viewing of a plurality of cameras. The system may be hardwired or wireless transmission may be utilized to further minimize the possibility of a malfunction at the onset of a catastrophic occurrence and to make the system more tamper resistant.\nThe comprehensive surveillance/communication of the subject invention supports communication of monitored data and/or commands or operational data between the ground or base station and the transport, between the transport and ground or terminal support vehicles and/or equipment, between the transport and various monitoring stations or systems, between transports, between the ground station and the support vehicles, between the monitoring station and support vehicles and between the monitoring stations or systems and the support vehicles. This permits the ground station to monitor and/or determine the identity, location, and heading of any vehicle in its range for tracking and collision avoidance, as well as monitoring sensor information, alarm conditions, emergency conditions, servicing requests, maintenance information, navigational information, requests for information such as flight plans, weather information, route maps, message traffic such as e-mail and the like. Similar information may be transmitted and received between transports, between transports and support vehicles and any of these and the ground station. The ground station may also send operational commands to the various monitoring systems both onboard the transport and ground mounted, such as camera tilt, pan and zoom and sensor activation. Other command signals such as xe2x80x9clock-onxe2x80x9d a specific condition or transport, sensor download, activation such as xe2x80x9clights-onxe2x80x9d or alarm (e.g., siren) activation and the like.\nIn a typical application, when an alarm from a specific transport is sent to the ground station it will be tagged with the G.P.S. coordinates of the transport. The alarm will also be reported to a security system, typically including a computerized center that distributes the information of the wireless LAN and where used, the wired LAN. The mobile and/or personal security units will also report their G.P.S. coordinates to the central computer so that the system knows the location of all security personnel at ant point in time. Once the alarm signal is received, the system can search and identify the closest appropriate personnel and alert them of the alarm condition. This is accomplished by calculating the length of the vectors between the transport G.P.S. and the various personnel G.P.S. signals. The shortest vectors are the nearest personnel and these can be alerted to respond to the alarm condition.\nThe selected personnel are then signaled by the security system of the present invention to respond. Audio, text and graphic communications may be utilized to inform the selected personnel of the condition and location. The system can also use its xe2x80x9cmappingxe2x80x9d function to assist the personnel in determining the best route to take in response. Because of the comprehensive nature of the system of the subject invention, both audio and image conditions of the transport can be communicated directly to the selected personnel, using video conferencing compression techniques of the LAN. If the desired, the personnel can switch cameras to obtain different views, or gain control of the steerable camera disclosed herein and survey the scene as appropriate via remote control. The two-way communication capability of the system would also permit the personnel to communicate conditions and the need for additional personnel or equipment both to the system computer and directly to other personnel.\nThe security computer system will register the G.P.S. location of the selected personnel as well as the location of additional or xe2x80x9cback-upxe2x80x9d personnel in order to coordinate their movements and actions. The system can then provide essential audio, video and communications to the selected back-up personnel in order to coordinate the entire operation. The coordinates of fixed sensors may also be entered into the system so that personnel can determine the proximity of each available sensor to his G.P.S. location.\nIt should be noted that the request for back-up can be programmed to be automatically activated under certain conditions. For example, if a security personnel personal system detects an explosion or a gunshot, an automatic alarm condition can be activated to alert central security other personnel in the vicinity to indicated xe2x80x9cpotential bomb blastxe2x80x9d or xe2x80x9cpotential automatic weaponxe2x80x9d, all based on interpretation of the audio signal which is picked up by the sensors by comparing them to known acoustic signatures of these, types of events. Also, audio detectors or multiple detectors can precisely determine the precise instant in time of an audio event by use of synchronized real time clocks associated with each sensor. The precise timing of the event can then be transmitted among the sensors or to a centralized monitoring computer where the exact location of the source of the audio event can be determined through triangulation techniques. This exact location can then be transmitted and/or displayed on various monitor devices to direct security and emergency personnel to respond.\nShock and vibration detectors may also be included both on board, at fixed locations on the ground and in the portable or mobile units. For example, if a personal unit is dropped, an alarm would be generated. Smoke and heat detectors may also be incorporated to monitor the safety environment of personnel.\nIt is a primary object of the subject invention to provide the means and apparatus utilizing IP protocols for a comprehensive, multi-media, wireless surveillance and monitoring system for monitoring and tracking a commercial transport vehicle while in port or while in route.\nIt is a further object and feature of this invention to provide a comprehensive surveillance and monitoring system utilizing IP protocols supported by a wireless transmission system whereby communication of all data including live video and/or audio transmissions can transmitted between the transport, ground or base stations, remote sensor systems, remote or mobile monitoring systems and other transports.\nIt is another object and feature of this invention to provide tracking capability utilizing IP protocols to assure that a transport stays in an assigned zone while either in route or in the port or terminal.\nIt is a further object and feature of this invention to provide communication capability utilizing IP protocols for monitoring and/or responding to supply needs on board the transport in order to permit support personnel to expedite response and/or re-supply when the transport arrives in port.\nIt is also an object and feature of this invention to provide for monitoring of situational conditions of and surrounding the transport both while in port and while in route.\nIt is yet another object and feature of this invention to provide means utilizing IP protocols for archiving performance parameters for later recall in order to review performance and/or reconstruct events.\nIt is an additional object and feature of this invention to provide a ground surveillance and security system utilizing IP protocols for detecting the breach of commercial transport security while the commercial transport is on the ground or in a port or terminal and is unattended.\nIt is another object and feature of the subject invention to identify that a commercial transport is on the ground and needs to be monitored for tracking its exact location, and its orientation on the ramp.\nIt is also an object and feature of the subject invention to provide a security system utilizing IP protocols, which is integral with the commercial transport for providing ground security.\nIt is a further object and feature of the subject invention to provide communications utilizing IP protocols between the commercial transport and a ground security station to assure commercial transport security while the commercial transport is parked or unattended.\nIt is another object and feature of the subject invention to provide a comprehensive system utilizing IP protocols, multi-media data generating, collecting, displaying, transmitting, receiving and/or storage safety and/or surveillance scheme for commercial transport.\nIt is also an object and feature of the subject invention to provide an on ground security system utilizing IP protocols which incorporates the in-flight surveillance system in order to minimize the number of additional components required to implement the system.\nIt is also an object and feature of the subject invention to store video, images, audio and/or transducer data on the commercial transport being protected and/or at the ground security station.\nIt is yet another object and feature of the subject invention to provide apparatus utilizing IP protocols for permitting ground and/or base personnel to receive video, images, audio information and/or data relating to critical components and areas of a commercial transport and operational data such as dispatch information.\nIt is still another object and feature of the invention to permit the monitoring, storing and retrieval of any of a variety of video, images, audio signals and/or performance data by the tracking, surveillance and/or imaging equipment on board the commercial transport.\nOther objects and features of the subject invention will be readily apparent from the accompanying drawings and detailed description of the preferred embodiments."}
{"text": "In recent years, implementation of “cloud-based” services, high-performance computing (HPC) and other activities employing data centers and the like have seen widespread adoption. Under a typical data center installation, a large number of servers installed in server chasses and server racks are interconnected in communication using network links (e.g., Ethernet or InfiniBand) and various switching mechanisms, such as switch blades/modules and “top-of-rack” (ToR) switches.\nUnder aspects of HPC, a very large number of compute nodes are implemented to solve various tasks in a parallel or substantially parallel manner. In essence, each compute node performs one part of a larger, more complex task. In order to implement this scheme, there needs to be input data to and output data must be exchanged among compute nodes. This data is communicated using one or more interconnects.\nVarious types of interconnects are used to interconnect the computer nodes in an interconnect hierarchy. For example, at the top of the hierarchy are interconnects between computing cores in the same processor. At the next level are interconnects between processors in the same platform, such as server blade or module. Next are interconnects between platforms, such as a backplane in a blade server. This is followed by interconnects between server chassis and/or ToR switches, interconnects between server racks, and finally interconnects between data centers. Generally, the communication bandwidth between nodes is reduced as one moves farther down the hierarchy.\nIn addition to latencies corresponding to transfers across the interconnect links themselves (which is a function of the link bandwidth and length and switching speed), significant latencies result from operations performed at the interfaces to the interconnects and/or additional processing/required to prepare the data for transfer across various interconnect links in the interconnect hierarchy. These data transfer latencies collectively reduce the communication performance and therefore the performance of the overall HPC implementation, and may represent a significant portion of the total latency (processing and data transfer) for the compute nodes.\nAnother important aspect of HPC is the software architecture. Implementing software using 10's of thousands of compute nodes in a parallel manner requires a significantly different software architecture than that employed for conventional applications. In addition, specific software modules have been developed for using corresponding types of interconnects, such as software modules used for communicating over InfiniBand."}
{"text": "The disclosures herein relate generally to devices for aiding human sleep, and more particularly, to sleep systems that provide improved support for the human head and neck to promote better sleep."}
{"text": "1. Field of the Invention\nThe present invention is directed to a dual-belt shuttle unit, more specifically, a horizontal dual-belt shuttle unit extending in a horizontal plane in a production line setting. The dual-belt shuttle unit typically transfers an article lifting device in a horizontal path from a first work station to a second work station.\n2. Description of the Prior Art\nDual-belt shuttle units of the general type with which the present invention is concerned typically include a fixed frame extending between a first and second work station and a single carriage dual-belt shuttle unit for transporting an article lifting device between the work stations. A drive unit is permanently mounted to the fixed frame and drives the conveyor system which transfers the article lifting device in a horizontal plane on the carriage dual-belt shuttle unit.\nA disadvantage of this previous system is that it requires a separate fixed frame and carriage dual-belt shuttle unit, along with a continuous support frame structure spanning several work stations. Further, each drive unit is permanently disposed on the fixed frame member thereby eliminating any possibility of exchanging a heavy duty drive unit for a lightweight drive unit when lifting and transporting different sized workpieces.\nThe present invention is directed to a dual-belt shuttle unit for horizontally transferring an article lifting device able to support multiple carriage dual-belt shuttle units for transferring the article lifting device between multiple work stations. The fixed frame member of the dual-belt shuttle unit is extendable along a horizontal plane between each work station. If necessary, the frame members may be attached in series along the work path to create a single assembly line as necessary. Each unit is self-supporting and, therefore, easily replaceable. Furthermore, each unit contains its own separate drive unit which itself is demountable and readily replaceable."}
{"text": "1. Field of the Invention\nThe present invention relates to a substrate processing apparatus that performs processing on a substrate, a storage device and a method of transporting a substrate storing container.\n2. Description of the Background Art\nSubstrate processing apparatuses are used to subject various substrates such as semiconductor substrates, substrates for liquid crystal displays, plasma displays, optical disks, magnetic disks, magneto-optical disks, photomasks, and other substrates to various types of processing.\nA substrate storing container is used for carrying a plurality of substrates in/out of the substrate processing apparatus from/to the outside. Examples of such a substrate storing container include a FOUP (Front Opening Unified Pod). The substrate storing container is provided with an opening through which the substrates are taken out or stored, and a cover for opening/closing the opening.\nAn opener places the substrate storing container thereon and opens/closes the cover of the placed substrate storing container. In the substrate processing apparatus including the opener, the substrate storing container storing the substrates before processing is placed on the opener of the substrate processing apparatus, for example. In this state, the cover of the substrate storing container is opened, and the substrates before processing in the substrate storing container are taken out to the substrate processing apparatus. This causes the taken out substrates to be subjected to given processing. Then, the substrates after the processing by the substrate processing apparatus are stored in the empty substrate storing container. With the cover closed, the substrate storing container storing the substrates after the processing is transported from the opener to the outside of the substrate processing apparatus (JP 2003-257945 A, for example).\nIn recent years, processing capability of the substrate processing apparatus has been improved with an increasing number of stacked processing units that perform processing on substrates. This allows more substrates to be processed by the substrate processing apparatus per unit time, thus requiring more substrates to be taken out from the substrate storing container to the substrate processing apparatus. The number of substrates to be stored from the substrate processing apparatus to the substrate storing container per unit time also needs to be increased. Therefore, the number of openers on which the substrate storing container is to be placed is preferably increased according to improved processing capability of the substrate processing apparatus.\nHowever, the limited number of openers can be mounted because of footprint constraints. Therefore, the throughput of the substrate processing apparatus may be limited due to exchange time of substrate storing containers."}
{"text": "In a number of industries, it is important to determine and map out-of-plane displacements of structural panels during loading. This is particularly important in the aerospace industry where it is important to determine displacements of load-sensitive structures such as composite panels.\nTraditionally, the methods used to predict buckling failures of composite structures have been largely empirical. The predictions were based on finite element analysis and a large number of failing load tests performed on stiffened shear panels, compression panels, and beams in universal test machines. The tests provide very little information on the deformation of the structures as a function of applied load magnitude. Analysts and designers have virtually no test correlation of any parameter except failing load.\nThe prior art has suggested the utilization of the projection Moire mapping technique for determining out-of-plane displacements of structural surfaces. Typically, this is done by projecting a line ruling on the surface of the structure to be examined. The surface should have a color and texture which permits the line ruling to be imaged on it. Ideally, the surface should be painted with flat white paint. The image of the ruling is viewed by a lens of a recording station and imaged on another ruling located at a recording station image plane. If the image of the projected grille on the recording station image plane and the recording station master grille have different pitches, Moire fringes are produced. If they are of the same pitch, no fringes will be produced at the recording station. In the latter situation as a test surface is deformed away from its initial position, Moire fringes proportional to out-of-plane displacement will be created at the recording station image plane. The resulting recorded image will be a contour map of the deformed surface and sequential photographs may be taken at the recording station image plane as the load is varied on the structural member being observed. The information derived from the resulting maps are useful in analyzing the effects of load, particularly shear loads which tend to cause buckling in composite structural elements."}
{"text": "A number of viruses have been reported to cause lethal hemorrhagic disease in humans and certain other primates. These viruses include viruses from a number of viral families, which include Filoviridae, Arenaviridae, Bunyaviridae, and Flaviridae. In particular the Filoviridae, which include the Ebola and Marburg viruses, have resulted in significant morbidity and mortality in infected populations. At present there are no satisfactory disease specific therapies or vaccines for diseases caused by the Filoviridae.\nPatients affected with hemorrhagic fevers develop a severe consumptive disseminated intravascular coagulation.\nDisseminated intravascular coagulation (DIC) is typically characterized by wide-spread systematic activation of the coagulation cascade resulting in excess thrombin generation. In addition, activation of the fibrinolytic system coupled with the consumption of coagulation factors results in a depletion of clotting factors and degradation of platelet membrane glycoproteins.\nConventional treatment of DIC has been aimed primarily at treatment of the underlying etiologic disease process and secondarily at the coagulopathy that results in the thrombotic and hemorrhagic manifestations. Reported therapies include replacement therapy of coagulation factors by transfusion of fresh frozen plasma. Heparin has been reported as sometimes used in combination with replacement therapy.\nTissue factor (“TF”) is a 47 kDa transmembrane glycoprotein that is the major cellular trigger of blood coagulation under physiologic conditions. The factor VIIa-tissue factor (“fVII/TF”) catalytic complex is able to generate factor Xa via direct activation of factor X, and indirectly through the activation of factor IX, thus initiating thrombin generation. It has been reported that tissue factor also plays an important role in disease processes resulting from the activation of the coagulation pathway. For example, TF levels are reported to be elevated during bacterial sepsis and this is believed to contribute directly to the pathogenesis of multiple organ failure (Doshi et al., “Evolving role of tissue factor and its pathway inhibitor”, Crit. Care. Med. (2002) 30(5 Suppl): 5241–5250). In addition, a number of viruses have been reported to activate the coagulation system following infection; such activation may also be triggered by the up regulation of TF expression (Bowman, et al., “Procoagulant and inflammatory response of virus-infected monocytes”, European J. Clin. Invest. (2002) 32: 759–766; Baugh et al., “Regulation of extrinsic pathway factor Xa formation by tissue factor pathway inhibitor”, J. Biol. Chem. (1998) 273: 4378–4386; Taylor et al., “Active site inhibited factor VIIa (DEGR VIIa) attenuates the coagulant and interleukin-6 and -8, but not tumor necrosis factor, responses of the baboon to LD100 Escherichra coli”, Blood (1998) 91: 1609–1615). A variety of inflammatory stimuli, including bacterial cell products, viral infection and cytokines, have been reported to induce the expression of TF on the surface of endothelial cells and monocytes, thereby activating the coagulation pathway (Doshi, et al.).\nTissue Factor Pathway Inhibitor (TFPI) is an endogenous systemically circulating plasma protein which is said to function as a physiological anticoagulant by inhibiting VIIa/TF complexes and preventing the initiation of coagulation. TFPI contains multiple Kunitz-type protease inhibitor domains, and is the principal physiologic inhibitor of TF/FVIIa. TFPI has been reported to bind to and inhibit factor Xa directly, prior to forming a quaternary inhibitory complex with TF/FVIIa, thereby inhibiting thrombin generation (Baugh et al.). Processes for preparing recombinant TFPI have been reported. See, e.g. U.S. Pat. No. 6,300,100 to Kamel et al.\nIn addition to its role in initiating coagulation, the TF/FVIIa has been reported to have direct proinflammatory effects independent of the activation of coagulation in man (Taylor et al.). In experimental settings where animals were depleted of TFPI, the animals were reported to have a demonstrated sensitivity to bacterial endotoxin and a higher propensity to develop intravascular coagulation (Doshi et al.). In a lethal E. coli sepsis model in baboons, treatment with TFPI was reported to attenuate the procoagulant and inflammatory cytokine interleukin-6 (IL-6) response and to prevent mortality (Creasey et al., “Tissue factor pathway inhibitor reduces mortality from Escherichia coli septic shock”, J. Clin. Invest. (1993) 91(6): 2850–2860). However, in healthy human volunteers administered low dose bacterial endotoxin, blocking TF/FVIIa with TFPI was reported to have no impact on inflammatory cytokines, but to completely prevent endotoxin-induced activation of coagulation (deJonge et al., “Tissue factor pathway inhibitor does not influence inflammatory pathways during human endotoxemia”, J. Infect. Dis. (2001) 183(12): 1815–1818). A phase II clinical trial of recombinant TFPI (rTFPI) in patients with severe sepsis, the rTFPI group was reported to demonstrate accelerated decrease of IL-6 plasma levels (Reinhart et al., “Assessment of the safety of recombinant tissue factor pathway inhibitor in patients with severe sepsis, a multicenter randomized, placebo controlled single blind, dose escalation study”, Crit. Clin. Med. (2001) 29(11): 2081–2089). A phase III trial of trifacogin, a rTFPI, in severe sepsis failed to show a reduction in the primary end-point of 28-day all cause mortality (Doshi et al.).\nProtein C (PC) is another component of the natural anticoagulant system in mammals. Unlike TFPI, which acts at the level of TF/fVIIa, the PC pathway is reported to inhibit coagulation by down regulating thrombin formation via the proteolytic inactivation of the non-enzymatic cofactors factor Va and factor VIIIa (Esmon C., “Protein C pathway in sepsis”, Annals of Medicine (2002) 34:598–605). The active component of the PC pathway is termed activated PC (aPC). aPC is formed by the action of thrombin bound to the non-enzymatic co-factor thrombomodulin. A recombinant form of aPC (drotrecogin alfa) has been reported to reduce the incidence death in patients suffering from severe bacterial sepsis (Bernard et al., “Efficacy and safety of recombinant human activated Protein C for severe sepsis”, N. Engl. J. Med. (2001) 344:699–709).\nPublished United States Patent Application, publication number US 2001/0028880 A1, is said to relate to the treatment of viral hemorrhagic fever with Protein C."}
{"text": "The present invention relates to rechargeable appliances and more particularly but not exclusively to an electric tooth brush and charger therefor.\nTypically, rechargeable tooth brushes are provided with a stand that rests on a bench surface. The stand has a primary coil that interacts with a secondary coil in the tooth brush to charge the batteries within the tooth brush. Extending from the stand is an electric cord terminating with the plug that is engaged within a wall socket to provide the primary coil in the stand with alternating electric power.\nThe above-described charger for tooth brushes suffers from the disadvantage that the stand occupies bench space and inhibits cleaning of the bench space. This problem is exacerbated by the cord.\nAlso known are wall-mounted chargers. These have the disadvantage that attachment to a wall is required. They also require a cord.\nIt is an object of the present invention to overcome or substantially ameliorate the above disadvantages.\nThere is disclosed herein a battery charger for an article having a charging secondary circuit including a secondary coil, the charger including:\na body having a socket with a passage to receive the articles so that the article is supported by the socket with the secondary coil located at least partly in the socket;\na charging primary circuit in the body and including a primary coil, the primary coil being positioned so as to surround the passage so as to be adjacent the secondary coil when located in the socket;\nat least two connection pins rigidly mounted in the body to engage a wall-mounted AC power supply, the pins being electrically connected to the primary circuit; and wherein:\nthe body and pins are adapted to engage the power supply so that the body is supported thereby together with the article supported in the socket.\nPreferably, the socket includes an open top and open bottom to receive the article so that the article is generally vertically oriented.\nPreferably, the body further includes a sleeve mounted in the socket, the sleeve having an open top and open bottom to receive the article so that the article is generally vertically oriented.\nThere is further disclosed herein a combination of the above battery charger and the article, and wherein the article is a rechargeable tooth brush including rechargeable batteries electrically associated with the secondary circuit."}
{"text": "The present invention relates to substantially anhydrous binders based on polyisocyanates which are modified by the addition of an organic sulfonic acid such that they have a self-releasing effect. The invention is also directed to the use of such binders in the production of shaped articles such as panels by the hot pressing of a wide variety of organic and/or inorganic materials such as substances containing lignocellulose.\nPressed materials such as chip boards, composite panels and other shaped articles are usually produced by hot pressing the inorganic or organic raw material such as a composition of wood shavings, wood fibers or another material containing lignocellulose, with binders such as aqueous dispersions or solutions of urea/formaldehyde or phenol/formaldehyde resins. It is also known to use isocyanates or isocyanate solutions as binders for pressed panels instead of formaldehyde resins (German Auslegeschrift No. 1,271,984; German Offenlegungsschriften Nos. 1,492,507; 1,653,177 and 2,109,686). The use of polyisocyanates as binders improves the stability and the moisture-resistance of the products and improves their mechanical properties. In addition, polyisocyanates have a wide range of processing advantages as binders as disclosed in German Offenlegungsschrift No. 2,109,686.\nThe large scale production of materials which are bonded with polyisocyanates, in particular materials containing lignocellulose such as chip boards, fiber boards or plywood is, however, impaired by the marked tackiness of the polyisocyanates. This tackiness, after the hot pressing treatment, causes the molded articles to adhere strongly to metal articles, in particular the steel or aluminum surfaces of the press. Such adherence thus makes it more difficult to remove the molded articles from the mold.\nPreviously proposed methods of solving this mold release problem suffer from significant disadvantages. Release agents which have been developed especially for isocyanates frequently have a good release action but, in industrial applications, they are neither reliable nor economical enough and, additionally, may cause faulty bonding or difficulties in coating during the subsequent processing of the plates.\nIt has been proposed in German Offenlegungsschrift No. 1,653,178 that, during the production of panels or shaped articles by the hot pressing of mixtures of materials containing lignocellulose and polyisocyanates, the surfaces of the press or pressing molds be treated with polyhydroxyl compounds such as glycerin, ethylene glycol or polyester and polyether polyols before the pressing operation. A disadvantage of this process is that a separate operation is required to apply this release agent and, in addition, a proportion of the polyisocyanate is consumed by the reaction with the release agent. According to German Offenlegungsschrift No. 2,325,926, another possible way of improving the release behavior of the shaped articles involves using as release agents those compounds which act as catalyst with isocyanates to form isocyanate. However, a disadvantage of this approach is that the catalysts have a destabilizing effect on the isocyanate and thus prevent the formation of a suitable isocyanate binder.\nAn object of the present invention is to overcome the above-mentioned disadvantages in the production of shaped articles using polyisocyanates by providing binders based on isocyanates which may be stored and which ensure that the pressed articles may be removed from the mold without difficulty. It has surprisingly been found that this object may be achieved if the isocyanate is used in combination with an organic sulfonic acid."}
{"text": "The present invention relates to a surface neuroprosthetic device for muscle stimulation.\nMovement impairment in a limb can result from a number of neurological or orthopedic pathological conditions. Activation of selected muscles of the limb by functional electrical stimulation (FES) to generate controlled movements has been used, both as a therapeutic modality and for the improvement or restoration of activities of daily living (ADL) or functional restoration.\nDevices based on surface electrical stimulation which have been developed for activating specific body sites include the dropfoot system, which activates the ankle joint, modifying hemiplegic gait; hybrid FES-orthosis systems for restoring gait in spinal cord-injured patients, and systems for therapeutic activation and functional restoration of the hand, such as the Handmaster, described in U.S. Pat. No. 5,330,516 (Nathan). Other prior art devices are described in U.S. Pat. No. 5,562,707 (Prochazka); U.S. Pat. Nos. 4,580,569 and 4,558,704 (both to Petrofsky). Use of these devices requires them to be placed on the body limb and to be removed, possibly once or several times, each day. One of the main barriers to widespread use of surface neuroprostheses has been the lengthy time and high level of expertise required to position the electrode array and to hold the electrodes on the limb surface. This is generally carried out by the patient himself, who suffers from neurological or orthopedic deficits. A hemiplegic patient, having the use of only one hand, will often be required to don and doff the neuroprosthesis by himself. Quadriplegic users may be able to don or doff the neuroprosthesis, or may require assistance.\nThe soft tissue of the paralyzed or plegic limb is highly flexible and mobile. Mechanical interaction between the limb and the device during the donning of the neuroprosthesis can result in geometric distortion of the limb's soft tissue. This problem should be addressed, as neuroprostheses are intended to be worn for several hours each day."}
{"text": "The present disclosure relates generally to retrieving and storing information pertaining to parts in a vehicle, such as an aircraft, and more specifically to retrieving and storing data based on locations of aircraft parts in one or more three dimensional models.\nAircraft are modeled and designed using three dimensional (“3D”) computer aided design (“CAD”) software. As a result, virtually every part of the aircraft is defined in 3D, including structure, mechanical systems, electrical systems, and propulsion systems. Within a 3D CAD database, each part may have various attributes associated with it, such as geometry information (e.g., centroid coordinates, shape, and volume), weight information, and vendor information. Many 3D CAD database systems enable a user to use 3D navigation to find data stored within the CAD database that is associated with a part or parts. However, it is not always convenient to store all information associated with aircraft parts in a 3D CAD database or in a format that indexes information in association with 3D locations within an aircraft. Accordingly, aircraft manufacturers (OEM) and/or airlines who own and/or operate aircraft may use multiple database systems that have different formats and that are indexed in different manners to store and deliver information about aircraft parts (e.g., manufacturing data, aircraft maintenance data, aircraft repair data, and/or spare parts data). As a result, aircraft part information stored in one database, such as a maintenance database, may not be accessible to a user using a system interface designed for use with another database, such as a 3D CAD database."}
{"text": "In network computing, it is common to discover and catalog software or hardware services that one node on a network might offer to other nodes on the network, such as file conversion services, data streaming services, database services, printing services and the like. These services ordinarily must be synchronized across the nodes, so as to ensure that each node has a catalog of services that is the same as all other nodes on the network.\nIn one example, each node tracks local changes by incrementing a counter. The node determines that synchronization with another node is necessary if the local counter does not match a counter at the other node."}
{"text": "The practical limits of optical lithography assume that the medium through which imaging is occurring is air. This practical limit is defined by the effective wavelength equation\n            Λ      eff        =          λ              2        ·        n        ·        NA              ,where λ is the wavelength of incident light, NA is the numerical aperture of the projection optical system, and n is the index of refraction of the medium. Now, by introducing a liquid (instead of the air) between a last lens element of the projection optical system and a wafer being imaged, the refractive index changes (increases), thereby enabling enhanced resolution by lowering the effective wavelength of the light source. Lowering a light source's wavelength automatically enables finer resolution of smaller details. In this way, immersion lithography becomes attractive by, for instance, effectively lowering a 157 nm light source to a 115 nm wavelength, thereby gaining resolution while enabling the printing of critical layers with the same photolithographic tools that the industry is accustomed to using today.\nSimilarly, immersion lithography can push 193 nm lithography down to 145 nm. In theory, older technology such as the 193 nm tools can now still be used. Also, in theory, many difficulties of 157 nm lithography—large amounts of CaF2, hard pellicles, a nitrogen purge, etc.—can be avoided.\nHowever, despite the promise of immersion photolithography, a number of problems remain, which have so far precluded commercialization of immersion photolithographic systems. These problems include optical distortions. For example, during immersion lithography scanning, sufficient g-loads are created that can interfere with system performance. These accelerative loads can cause a vibrational, fluidic shearing interaction with the lens resulting in optical degradation. The up and down scanning motions within the lens-fluid environment of Immersion Lithography can generate varying fluidic shear forces on the optics. This can cause lens vibrational instability, which may lead to optical “fading”. Other velocity profile non-uniformities can also cause optical distortions."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an oil activated fuel injector and, more particularly, to an oil activated electronically or mechanically controlled fuel injector control valve which substantially eliminates captured air within working fluid of the fuel injector.\n2. Background Description\nThere are many types of fuel injectors designed to inject fuel into a combustion chamber of an engine. For example, fuel injectors may be mechanically, electrically or hydraulically controlled in order to inject fuel into the combustion chamber of the engine. In the hydraulically actuated systems, a control valve body may be provided with two, three or four way valve systems, each having grooves or orifices which allow fluid communication between working ports, high pressure ports and venting ports of the control valve body of the fuel injector and the inlet area. The working fluid is typically engine oil or other types of suitable hydraulic fluid which is capable of providing a pressure within the fuel injector in order to begin the process of injecting fuel into the combustion chamber.\nIt has been found in open systems that air becomes captured and locked within the grooves or orifices of the control valve (and a spool) during the venting of the working fluid during and at an end of a fuel injection cycle. This is mainly due to the fact that vent holes which surround the control valve body allow air to enter the system. This air will mix with the working fluid during the fuel injection process resulting in variations in fuel injection quantities. Of course, this will lead to inefficient shot to shot variations.\nBeing more specific, a driver will first deliver a current or voltage to an open side of an open coil solenoid. The magnetic force generated in the open coil solenoid will shift a spool into the open position so as to align grooves or orifices (hereinafter referred to as xe2x80x9cgroovesxe2x80x9d) of the control valve body and the spool. The alignment of the grooves permits the working fluid to flow into an intensifier chamber from an inlet portion of the control valve body (via working ports). The high pressure working fluid then acts on an intensifier piston to compress an intensifier spring and hence compress fuel located within a high pressure plunger chamber. As the pressure in the high pressure plunger chamber increases, the fuel pressure will begin to rise above a needle check valve opening pressure. At the prescribed fuel pressure level, the needle check valve will shift against the needle spring and open the injection holes in a nozzle tip. The fuel will then be injected into the combustion chamber of the engine.\nTo end the injection cycle, the driver will deliver a current or voltage to a closed side of a closed coil solenoid. The magnetic force generated in the closed coil solenoid will then shift the spool into the closed or start position which, in turn, will close the working ports of the control valve body. The working fluid pressure will then drop in the intensifier and high-pressure chamber such that the needle spring will shift the needle to the closed position. The nozzle tip, at this time, will close the injection holes and end the fuel injection process. At this stage, the working fluid is then vented from the fuel injector via vent holes surrounding the control valve body.\nReferring now to FIG. 1A, in current designs the vent holes 10 surround the control valve body 12 and the spool 14 such that air 16 in the control valve body 12 is below the working fluid level 18. This causes the grooves 20 of the control valve body 12 and the spool 14 to be filled with air 16. Now, during the next cycle time (as seen in FIG. 1B) when the spool 14 is shifted to the open position, this air 16 becomes locked within the grooves 20 causing air bubbles 22 to be formed within the working fluid 18 of the working ports 23. In order to inject fuel within the combustion chamber, this captured air will have to be compressed by the working fluid and dissolved partially into a dilution prior to the working fluid acting on the intensifier piston. This causes a shot to shot fuel variation (depending on the quantity of air in the working fluid) thus resulting in decreased fuel efficiency especially for low fuel quantities.\nThe present invention is directed to overcoming one or more of the problems as set forth above.\nIn a first aspect of the present invention, a check valve body has an inlet area and a working port in fluid communication with the inlet area. The working port is adapted to provide working fluid to an intensifier chamber of the fuel injector. At least one communication port is in fluid communication with the inlet area and the working port. At least one vent hole is provided which prevent air from mixing with the working fluid.\nIn another aspect of the present invention, the check valve body has an oil inlet area and a at least one port in fluid communication with the oil inlet area. The port transport oil between the oil inlet area and an intensifier chamber of the fuel injector. An aperture having at least one communication port provides a flow path for the oil between the ports and the oil inlet area. A spool is positioned within the aperture and includes at least one fluid path which are in alignment with the communication port of the aperture when the spool is in the first position. Vent ports vent the oil from the control valve body and prevent air from entering the at least one fluid path of the spool.\nIn still another aspect of the present invention, a fuel injector having a control body is provided. The control body has an inlet area, working ports, communication ports and fluid paths, a spool and at least one vent hole. The at least one vent hole is positioned above the working ports to reduce captured air in the working ports during a venting process. The fuel injector also includes an intensifier body and a spring loaded piston and plunger within a centrally located bore of the intensifier body. A high pressure fuel chamber is also formed in the intensifier body. A nozzle having a fuel bore is in fluid communication with the high pressure chamber, and a needle is positioned within the nozzle. A fuel chamber surrounds the needle."}
{"text": "1. Field\nThis disclosure is generally related to data synchronization. More specifically, this disclosure is related to service discovery by executing synchronization of collections between peers on a network using exact match names.\n2. Related Art\nIn many computing applications, it is often important for peers on a network to synchronize their respective collections of data. The proliferation of digital content creates a vast number of collections which require reconciliation. Content-centric network (CCN) architectures have been designed to facilitate accessing and processing such digital content. A CCN includes entities, or nodes, such as network clients, forwarders (e.g., routers), and content producers, which communicate with each other by sending “interest” packets for various content items and receiving “content object” packets in return. CCN interests and content objects are identified based on a unique name, which is typically a hierarchically structured variable length identifier (HSVLI) which comprises contiguous name components ordered from a most general level to a most specific level.\nA current CCN synchronization protocol uses a longest-prefix match method, where an interest in “/parc/events/” will match both “/parc/events/calendar.txt” and “/parc/events/conference.txt.” This current synchronization protocol is based on responses that extend the name of an interest. This synchronization protocol is described in the following applications:                U.S. patent application Ser. No. 13/720,736, entitled “DYNAMIC ROUTING PROTOCOLS USING DATABASE SYNCHRONIZATION,” by inventors Van L. Jacobson and Marc E. Mosko, filed 19 Dec. 2012 (hereinafter “U.S. patent application Ser. No. 13/720,736”); and        U.S. patent application Ser. No. 13/681,306, entitled “DATA TRANSPORT BY NAMED CONTENT SYNCHRONIZATION,” by inventors Van L. Jacobson and Marc E. Mosko, filed 19 Nov. 2012 (hereinafter “U.S. patent application Ser. No. 13/681,306”);the disclosures of which are herein incorporated by reference in their entirety.        \nAs CCN architectures evolve, there is a need to discover devices and services offered by those devices based on a synchronization protocol that allows the use of exact name match, rather than the current longest-prefix match. In such a new architecture, the current CCN synchronization protocol will not work."}
{"text": "1. Field of the Technology\nThe present technology relates to a video game program and a video game system and, more particularly, to a video game program and a video game system for playing a video game using an optical pointing device.\n2. Description of the Background Art\nThere are conventional video game systems using optical pointing devices. For example, Patent Document 1 (Japanese Laid-Open Patent Publication No. 8-71252) discloses a shooter video game device using a gun-shaped controller. The gun-shaped controller has a CCD camera for sensing light-emitting members provided, as marker objects, around the video screen. Thus, the video game device can detect the distance between the screen and the gun-shaped controller, the rotation of the gun-shaped controller, and the position thereof. Specifically, Patent Document 1 discloses a shooter video game device capable of calculating the position (the coordinates thereof) on the screen being pointed at by the gun-shaped controller.\nPatent Document 2 (Japanese Laid-Open Patent Publication No. 2002-81909) also discloses another shooter video game device employing substantially the same configuration as that of Patent Document 1 (see paragraph 0020, etc.). With this shooter video game device, a position on the screen being pointed at by a gun-shaped controller is calculated by using an image captured by image capturing means provided in the gun-shaped controller. In the shooter video game, the calculated position is used as the position at which a gunsight image of the gun-shaped controller is displayed. There are known video games in which the position on the screen being pointed at by a gun-shaped controller is calculated, and used the position at which a gunsight image is displayed.\nWith these video game devices, the position on the screen being pointed at by the gun-shaped controller can be calculated, and the gunsight image is displayed at the calculated position on the screen. Therefore, the player can easily know the direction in which the gun is being directed, i.e., the position on the screen being pointed at by the gun-shaped controller. When the markers (light-emitting members) are out of the image capturing area of the image capturing means provided in the gun-shaped controller, it means that the aiming point of the gun-shaped controller is outside the screen, in which case the gunsight image on the screen cannot be moved or controlled. If the image capturing means again starts sensing the markers, the gunsight image will be displayed at the position on the screen pointed at by the gun-shaped controller.\nA controller using an optical pointing device as described above has primarily been used for moving a gunsight image on the screen in a shooter video game. Recently, with an increasing variety of types of video games, such a controller may be used for controlling any other object in the game space, in addition to controlling a gunsight image on the screen. For example, there are at least two ways in which an object in a three-dimensional game space is moved by using a controller. In the first way, an object is placed at a position in a three-dimensional game space corresponding to a position on the screen pointed at by the controller, and the object is moved across a two-dimensional plane in the game space according to a two-dimensional movement of the position pointed at by the controller. With this technique, however, the object can only be moved across a two-dimensional plane in the game space, which is not at all different from moving a gunsight image in conventional shooter video games, except that the game space is three-dimensional.\nIn the second way, the player can move the controller up and down to change the Y coordinate of the object placed in a three-dimensional game space defined in an XYZ coordinate system, thereby moving the object up and down in the game space. The player can move the controller left and right to change the X coordinate of the object in the game space, thereby moving the object left and right in the game space. Moreover, the player can move the controller forward and backward to change the Z coordinate of the object in the game space, thereby moving the object forward and backward in the game space.\nIn the second way as described above, the object in the three-dimensional game space may not always be displayed at the position on the screen pointed at by, for example, the controller (hereinafter referred to as the “pointed position”, corresponding to the position of the gunsight image on the screen in conventional shooter video game devices). With a typical three-dimensional game space, the game space is shown in a bird's-eye view as viewed from a virtual camera. Therefore, the up/down, left/right and forward/backward movements in the game space do not always coincide with those as viewed from the virtual camera. Specifically, when an object moves in the up/down direction by a certain distance in the game space, the object moves by a shorter distance on the bird's-eye view of the game space as viewed from the virtual camera. When the object moves in the depth direction in the game space, which is the direction away from the virtual camera, the object on the screen as viewed by the virtual camera moves toward the center of the screen. Thus, the position on the screen pointed at by the controller may not coincide with the position at which the object is displayed on the screen.\nIn conventional shooter video game devices, the gunsight image is displayed at the position pointed at by the controller. Therefore, the player can easily recognize the direction in which the controller is being directed. However, where the controlled object is not displayed at the pointed position, the player cannot recognize the direction in which the controller is being directed. Since the player cannot recognize the pointed position by looking at the screen, the player may move the controller to such a position that the image capturing means of the controller can no longer sense the markers, whereby the object can no longer be controlled. Particularly, if the player cannot identify the current direction of the controller (i.e., the pointed position) when starting to control the object, the player will have to control the object without knowing at all the direction in which the controller is being directed. Therefore, it is likely that the player moves the controller to such a position that the image capturing means of the controller can no longer sense the markers. This also significantly degrades the controllability in using the controller.\nAs described above, if the display position of an object being controlled by using a controller does not coincide with the position pointed at by the controller, the player may move the controller to such a position that the object can no longer be controlled."}
{"text": "This application is a 371 of PCT/KR01/01059 filed Jun. 21, 2001.\nThe present invention relates to 3,9-diferulylcoumestrol and a skin cosmetic product containing the same, more particularly to a novel skin cosmetic product capable of improving anti-aging, whitening and hair-seedling effect and remedying skin disorder from a cosmetic point of view.\nSkin aging means a condition that wrinkles or droops are formed with age. The latest study about skin aging is largely divided into two ways; a histological approach studying for example the mechanism of matrix formation or destruction of connective tissue such as skin collagen and a immunological approach studying for example the relation between oxidizing stress and aging.\nAt first, the cause of skin wrinkle and droop in a histological viewpoint, is degeneration of collagenous fiber or ellastin fiber[M. Oshima et al., Matrix, 13, 187(1993); S. Harumiya et al., J. Biocham., 120, 745(1996)]. Further, the lowered supporting force of hypodermic tissue and muscle tissue increases wrinkle. In particular, collagen, which is the original material or collageneous fiber and ellastin, is synthesized in fibroblast and forms collagenous fiber and ellastin.\nCell membrane of fibroblast has signal receptor glycoprotein capable of specifically recognizing nutrition passage and extracellular materials. The receptor is connected to ion transport passage and if the exterior materials approaches the receptor of the fibroblast cell membrane, the signal is transported into the cell, finally to the inside of the nucleus. It is ordered to synthesize collagenous protein and more proteins than ordinary state are synthesized and released out of the nucleus. The protein grows by passing through the various organs of the cell and released from the cell. The released protein becomes a base of collagenous fiber or ellastin among cells and then elastic skin with a strong supporting power. [A. Hayasho et al., Fragnance J., 2, 32(1992)]. However, collagen biosysthesis in corium sharply decreases with age and various stress such as ultraviolet rays or environment pollution, resulting in forming wrinkle.[M.Oshima et al., Matrix, 13, 187(1993)].\nIn immunological viewpoint, oxidative stress in or out of human body destroys cells and biological connective tissue, which increases the number of wrinkles and makes them deeper. Furthermore, as the bio-synthetic ability of the biological anti-oxidant materials becomes lowered with age, the activation of oxidizing materials such as active oxygen which is introduced out of the body or metabolically inevitably generated cannot be prevented, resulting in accelerating wrinkle formation.[L. H. Kligman et al., J. Invest. Dermatol. 93, 210 (1989); H. Tanaka et al., Arch. Dermatol. Res., 285, 352(1993); M. D. Carbonare et al., J. Photochem. Photobiol. B. Biol., 14, 105(1992); T. Okada et al., Fragnance J., 11,23 (1993)]. Therefore, to prevent oxidization of body including skin, artificial intake or application of anti-oxidant agent is essential.\nThere are a number of materials with anti-oxidant function in nature. The principle of anti-oxidant function could be explained in the following way. Lipids forming cell membranes are easily attacked by peroxides or other active oxygen species generated from environmental factor such as ultraviolet rays or pollution or internal factor such mental stress. They are changed into peroxide lipid and forms radical polymer, resulting in destruction of cell membrane [T. Shibamoto et al., American Chem Soc 154(1994)]. The destruction of cell membrane protecting intracellular materials causes cell necrosis and skin aging. Anti-oxidant materials is oxidized instead before the cell membrane lipid is oxidized, and prevent lipids from being oxidized. So far, anti-oxidant materials used in cosmetics are various plant extract such as tocopherol and ascorbic aid. However, there is a problem that they are very unstable or do not show enough effect.\nIn present, for the most part of the functional materials applied to anti-aging cosmetics, the anti-aging effect have been discussed only in the basis of in vitro test results and clinic test results, instead of proving histological improvement, because of technical limitation. That is, the anti-aging materials being developed as cosmetics these days are focused only to the in vitro immunological effect. Therefore, it is not believed that such materials bring a substantial anti-wrinkle effect.\nUnder the above circumstance, the inventors invented 3,9-diferurylcoumestrol synthesized from ferulic acid and coumestrol with superior anti-oxidant force, which can solve chemical unstability or lack of effect and has histological and immunological effect for restraining aging and wrinkle formation.\nIt is not easy to develop materials with a substantial anti-aging function. As described above, the sources and products marketed as an anti-aging concept were concentrated on anti-aging effect, but could not obtain desired effects in practice.\nThe anti-oxidant effect can protect against the attack of internally generated or externally introduced active oxygen species to some extent However, such function has not guaranteed the substantial anti-aging effect in a histological standpoint, such as visible elimination of wrinkles. Therefore, it has been required to develop materials or products which have immunological anti-aging function by anti-oxidant function and further can accelerate skin collagen generation, so that skin recovering and skin elasticity and wrinkle improvement performance are all excellent\nThe inventors tried to develop material which can intensifies immunological function as a basic direction for anti-aging and histologically improve skin wrinkles. Considering the above, we came to synthesize a new concept of anti-aging material from ferulic acid known as its excellent anti-oxidant ability and coumestrol known as its excellent effect for accelerating collagen synthesis, as precursors.\nFerulic acid exists in asafetida, pine resin, leaf of rice plant, and is known as its antibiotic and anti-cancer ability and mild anti-oxidant property.[M. G. Smart et al., Aust J. Plant Physical. 6, 485(1979)]. Further, it is known as having a thrombus coagulation preventing function and used as an antiseptic of food.[T. Tsukiya et al., Jpn. Lolai 7518621(1975), C. A.83, 7602v(1975)]. Ferulic acid was extracted from a plant resin and named by H. Hiasiwetz and his colleagues of Australia in 1866 [H. Hiasiwetz et al., Ann. 138, 61(1866)].\nThe active oxygen species derived from oxygen molecule are essentially generated in the metabolism and their concentration should be maintained in the human body. Too low concentration has a bad effect on the metabolism and too high concentration causes peroxidization of bio-molecules such as cell lipid membranes of the organs[A. M. Kligman et al., J. Am. Acad. Dermatol., 15,836(1986); E. J. Van Scott et al., Cutis., 43, 222(1989)].\nFerulic acid has an mild oxidizing force over various active oxygen species derived from oxygen molecule and an strong oxidizing force over oxidative transitional metals such as Iron ion or Copper ion.[M. G. Smart et al., Aust J. Plant Physiol. 6, 485(1979)]. It is further reported that the Ferulic acid is used for preventing auto-oxidization of natural oil such as linseed oil. Ferulic acid has been proved that it is an active material as an ant-oxidant material in the filed of cosmetic, medicine and food.\nHowever, Ferulic acid has only simple anti-oxidant function and is mainly used as an additive for preventing oxidization of effective components of products by transitional metal or other introduced active oxygen species[T. Tsukiya et al., Jpn. Kokai 751862191975], C. A. 83, 7602v(1975)].\nCoumestrol is generally found in leguminous plants and has mild anti-oxidant effect and is known to have anti-inflammation, anti-fungal and anti-virus action [Darbarwar, M., et al., J. Sci. Ind. Res., 35, 297(1982); Jeffrey, B. H. et al., Phytochemical Dictionary, Taylor and Francis, London, pp 418(1933)].\nAlthough its accurate mechanism is not revealed, it is known to accelerate production of connective tissue such as collagen in skin and internal organs and believed to have an anti-aging or anti-wrinkle effect [Bickoff, E. M., et al., J. Anim. Sci., 19,4(1960)]. However, the amount of coumestrol contained in a black soybean is extremely small and thus simple extract method does not lead to desired effect Therefore, the objective of the invention is to develop a new material with substantial anti-aging effect and skin-elasticity improving effect together with whitening effect in cosmetics by producing coumestrol in a mass-synthetic way and connecting it with a material with an excellent anti-oxidant effect such as ferulic acid.\nAlthough about three synthetic method of coumestrol has been known. It is not generalized to synthesize coumestrol. The inventors used one of the three methods [E. M. Bickoff, et al., J. Am. Chem. Soc., 80, 4381(1958)]. The used method is total synthesis consisting of total 7 stages. The present invention produces a material having strong anti-oxidant effect together with anti-aging and whitening effect of skin by accelerating the regeneration of skin connective tissue."}
{"text": "Consumers may use a financial transaction account, which may be associated with an account number and/or transaction instrument (e.g., charge card, credit card, debit card, gift card, etc.), as a form of payment or for identification in various transactions. However, sufficient inventory management systems and tracking methods do not exist for distribution of many financial transaction instruments.\nDistribution of financial transaction instruments may involve the distribution of physical objects directly to retail locations or through a distribution chain. Distribution of, for example, charge cards may take place through various third party distribution centers and distributors in a supply chain. Such systems are not generally configured for detailed tracking of inventory, and systems are generally not in place or sufficient to facilitate providing timely or useful information to the issuer of the transaction instrument.\nThese limitations on the ability to track inventory have historically not been a strong concern because the product being distributed and/or the method of distributing the product did not raise significant security, loss prevention, and re-claim issues. For example, distributed cards have typically been closed cards. “Closed cards” are cards that may be restricted to use in a particular store or within a particular chain of a stores. One example of a closed card is a pre-paid gift card that may only be purchased at, and only be accepted at, a clothing retailer, such as The Gap®. In contrast, “Open cards” are cards that are generally accepted at different merchants. Examples of open cards include American Express®, Visa®, and Discover®, MasterCard® cards, which may be used at many different retailers and other businesses.\nA relatively low level of concern often exists for loss control and inventory management in the distribution of transaction instruments that are closed system, valueless, and/or non-expiring cards in comparison to open system, pre-paid, and/or expiring transaction instruments. As such, the conventional wisdom has been that detailed tracking of the distribution of prepaid cards, charge cards, credit cards, and the like, may not be cost effective. By way of example, closed cards generally do not have expiration dates; and because of this, there typically has been no reason to reclaim the unsold cards after a period of time has elapsed.\nHowever, it is presently becoming desirable to distribute open financial transaction instruments (e.g., open debit/credit cards) in a publicly accessible manner. For example, it is desirable to hang an open card on a rack in high foot traffic areas of a store. Such a method of distributing may tend to increase the effectiveness of card distribution. However, open cards may have an expiration date associated with the card. Furthermore, the distribution of open cards may increase the need for of security measures.\nThus, there exists a need for systems and methods that facilitate determining the location of a financial transaction instrument (e.g., a card) and reclaiming the financial transaction instrument in advance of its expiration date. More generally, there exist a need for inventory management systems and methods that can track the distribution of open financial transaction instruments."}
{"text": "In order to obtain a mechanical component of steel having a surface layer portion enhanced in strength, a treatment is sometimes performed to provide the surface layer portion with a layer having a higher nitrogen concentration than the remaining region, i.e., a nitrogen enriched layer. For example, a nitriding treatment is performed. A conventional nitriding method for steel is representatively a gas soft nitriding treatment in which steel is heated in an atmosphere containing ammonia or a similar gas serving as a source of nitrogen to cause the nitrogen to penetrate into a surface layer portion of the steel. However, when steel containing chromium in a large amount such as at least 3.75 mass %, for example, is used to produce a mechanical component, the mechanical component has a surface layer portion having chemically stable oxide film. As such, when the mechanical component formed of steel containing chromium in the large amount undergoes the gas soft nitriding treatment the surface layer portion is not penetrated by nitrogen and the nitrogen enriched layer is not formed.\nTo address this, a plasma nitriding process is proposed as follows: an object of steel to be treated is placed in a vacuumed furnace and a gas containing a gas serving as a source of nitrogen is introduced into the furnace, and between the object and a member, such as a wall of the furnace, disposed to face the object a difference in potential is caused to cause glow discharge to cause nitrogen to penetrate into a surface layer portion of the steel configuring the object (see Japanese Patent Laying-open No. 2-57675 (Patent Document 1) for example). The plasma nitriding process is controlled for example as based on a spectral analysis of glow discharge, a density of a current flowing in the object, or the like, as proposed in Japanese Patent Laying-Open No. 7-118826 (Patent Document 2) and Japanese Patent Laying-Open No. 9-3646 (Patent Document 3). This allows a mechanical component formed of steel containing at least 3.75 mass % of chromium to have a surface layer portion provided with a nitrogen enriched layer.\nFollowing the recent improvement in performance and efficiency of a mechanical device employing a rolling bearing, high durability in a severe environment tends to be required to the rolling bearing. More specifically, a rolling bearing employed in a contaminated environment penetrated by hard foreign matter may be damaged in an early stage (in an operating time shorter than a calculated life of the bearing) due to gripping of the foreign matter. Furthermore, a rolling bearing rotating at high speed may have smearing even if the bearing is under relatively small load. Furthermore, a rolling bearing used in an insufficiently lubricated environment may cause seizure. When the rolling bearing is used in a high-temperature environment of a temperature exceeding 200° C., for example, hardness of components (bearing components) constituting the rolling bearing may be reduced, to reduce the durability of the rolling bearing.\nWhen the bearing components are made of steel, strength at a high temperature can be improved for improving the durability of the rolling bearing in the high-temperature environment by adding at least 3.75 mass % of chromium to the steel thereby improving tempering softening resistance of the steel. In order to improve the durability in the contaminated environment, a treatment of forming nitrogen-enriched layers having higher nitrogen concentrations than the remaining regions on surface layer portions of the bearing components by performing nitriding, for example, can be employed.\nIn bearing components made of steel having a high chromium content of at least 3.75 mass %, for example, chemically stable oxide films are formed on the surface layer portions. When ordinary nitriding is performed on these components, therefore, nitrogen does not penetrate into the surface layer portions thereof, and no nitrogen-enriched layers are formed. In relation to this, there is proposed a countermeasure of forming nitrogen-enriched layers by performing a plasma nitriding process, as described above.\nIn order to improve seizure resistance, there is proposed a countermeasure of dipping balls serving as the bearing components in an organic phosphorus compound for forming reaction layers on the surfaces thereof (refer to Japanese Patent Laying-Open No. 9-133130 (Patent Document 4), for example).\nPatent Document 1: Japanese Patent Laying-open No. 2-57675\nPatent Document 2: Japanese Patent Laying-open No. 7-118826\nPatent Document 3: Japanese Patent Laying-open No. 9-3646\nPatent Document 4: Japanese Patent Laying-open No. 9-133130"}
{"text": "The invention relates to the field of animal feed grains and a method of treating them so as to enable the animal to assimilate a greater quantity of the nutrients contained in the grain.\nThe literature discloses studies and tests which relate to \"reconstitution\" of feed grains. The term is defined in a Texas A & M University press release No. 2565, dated June 1968, as \"consisting of adding water to air-dry grain to increase its moisture content to 25-30 percent, followed by controlled fermentation in the absence of air.\"\nThe release describes feeding tests in 48 animals which were divided into 6 uniform groups. Two groups which served as \"controls\" were fed ground, dry sorghum grain. Two other groups were fed moist sorghum grain which was reconstituted in whole form and was ground just before feeding. The remaining two groups received sorghum grain which had been ground and reconstituted. The first and last mentioned groups had substantially the same daily gain whereas the groups fed with reconstituted whole grain decreased the feed required per unit of gain by 11%. The \"reconstitution\" process was stated to require at least 21 days before the grain could be fed, and also required air tight silos for commercial practice of the process.\nThe release also states that the exact changes which occur in grain during reconstitution are not well understood, but that the process \"may be\" similar to that which takes place during germination of the grain, which will be discussed further below.\nThe literature contains other, later experiments which are along the lines of those contained in Release No. 2565, and which generally corroborate the findings thereof.\nFor example, a report of tests by the North Florida Experiment Station, Quincy, Florida, dated Jan. 8, 1969, indicates that comparable results were attained by the same process using shelled corn instead of sorghum grain. Another report of tests, published at a \"Beef Cattle Short Course\" at Gainsville, Florida, dated May 1, 1969, indicates shelled corn should contain 22 to 30% moisture because with moisture levels below 22%, adequate fermentation and softening of the grain does not occur. It is noted that the later reports appear to define \"reconstitution\" as the addition of water to dry grain rather than this step AND the fermentation step combined, as defined by Release No. 2565. Since the latter definition appears to be more widely accepted, this specification uses the term to indicate the addition of water rather than both the addition of water and subsequent fermentation. The \"Short Course\" report also indicates that the fermentation process may be terminated after 10 days rather than 3 weeks, without material impairment of the benefits stated in the Release.\nThe Short Course report discusses a number of methods for \"reconstituting\" dry grain prior to fermentation: (1) Spraying water through nozzles onto grain passing through an auger which has its discharge end elevated; (2) Immersion in a flooded pit from which the grain is removed by an auger; (3) A mixer truck; (4) A \"reconstituter\" in which grain is sprayed with water as it passes through; (5) Flooded on truck with perforated hose, to prevent water channeling down through grain and draining out of truck body before grain is wet; and (6) Water applied in blower as grain is put in structure (for fermentation).\nApplicant's view of the state of the art is that all of the known methods for reconstituting grains, including the methods noted above, attempt to supply moisture to the interior of each kernel through the existing, naturally provided entrance channels, which occur in the germ area. The hull constitutes a natural \"raincoat\" which prevents absorption of water except in the germ area.\nIf, as suggested in Report No. 2565, the process of reconstitution is similar to that of germination, the first water absorbed would be taken up and held by the embryo, which produces an enzyme. The latter is conducted to the alourone layer which secretes a hydrolytic enzyme, which coacts with water in the starch matrix to convert or liquify the starch of the kernel. Thus, considerable time is required to absorb the required amount of water and to distribute it substantially uniformly throughout the interior of the kernel. As indicated above, the time period is on the order of 10 days, and since open air storage for such period could lead to deterioration or spoilage of the grain, air tight enclosure of the grain is necessary. Further, the air tight storage must be maintained during the entire feeding period to prevent uncontrolled fermentation.\nThe art cited during prosecution of the above identified parent application of which this application is a division, while teaching various types of methods and apparatus for treating feed grains, does not teach or suggest a method or apparatus wherein portions of the hulls of the kernels of feed grain are loosened from the interior portion of the kernels to cause the hulls to blister and thereafter opening the blisters to perforate the hulls."}
{"text": "1. Field of the Invention\nThe present invention relates to a content processing apparatus, a content processing method, and a content processing system, and more particularly to facilitating use of identical content on different information processing apparatuses.\n2. Description of the Related Art\nAlong with the popularization of the internet and advances in multimedia technology, information exchange over networks is developing to incorporate not only exchange of simple text data such as e-mail but also exchange of various content data such as sound data, image data, and movie data that includes both sound and video. These content data are exchanged among various types of terminals such as cell phones and laptop personal computers (PCs). To display content properly on various different terminals, a method to distribute content data converted to format corresponding to display capability of terminal that requested the content data has been proposed (e.g., JP-2007-11760-A.)\nMeanwhile, apparatuses such as projectors and electronic whiteboards also use content data in addition to cell phones and laptop PCs. Also, the way the content is used is changing, from not simply passively viewing movies, listening to audio, and browsing images to being able to edit documents and images. Yet the technology described in JP-2007-11760-A is limited to passively displaying content, and moreover such technology cannot apply to apparatuses that do not have editing functions in the first place.\nFor example, in displaying and editing content in a format by using an electronic whiteboard in a conference room, to move to another conference room and continue the meeting, an apparatus that supports displaying and editing content in that format is necessary in the destination conference room. That means users cannot continue the meeting under the same conditions as in the original conference room in case the apparatus in the destination conference room lacks the capability to edit information even if content is distributed in a format adapted to the apparatus at the destination conference room.\nIt should be noted that this kind of issue can be a problem not only distributing information but also using content in a format on a different information processing apparatus."}
{"text": "1. Field of the Invention\nThe present invention relates to a navigation system for detecting the position of a moving body to inform a user of the current position of the moving body, and a map display apparatus using the navigation system.\n2. Description of Related Art\nVarious types of navigation systems have been recently known, and these navigation systems are used while mounted on moving bodies such as vehicles, ships, etc. When such a navigation system is mounted on a moving body, it performs arithmetic and logic processing on information provided from various sensors to detect the current position of the moving body and display the current position on a screen of the navigation system.\nThis type of navigation system includes a position detector for detecting the absolute position of the moving body, a storage device for storing map data which include two-dimensional vector data obtained by projecting roads and points on the ground, such as structures, etc., onto a mesh-segmented plane using a universal transverse Mercator projection method, and character data associated with the vector data, an input device for receiving external instructions (commands), and a display device for reading out desired vector data from the mesh-segmented map stored in the storage device in accordance with an instruction input from the input device, to perform data conversion processing on the data, and then displaying the map on the display.\nThe data conversion processing includes shift conversion processing for changing the display position on a map, scale conversion (enlarging/reducing) processing for displaying a map on any scale, and rotational coordinate-transformation processing for changing a display direction on the map. With these processings, a plan view map is obtained by drawing the ground through an orthogonal projection in a vertical direction.\nAs described above, when a map is displayed on the screen, the conventional navigation system as described above adopts a plan-map display mode in which the ground is drawn in the vertical orthogonal projection style. Therefore, in order to simultaneously display two spots which are far away from each other, a reduction-scale of the map must be necessarily increased, and thus detailed information cannot be displayed on the screen.\nAn object of the present invention is to provide a bird\"\"s-eye view forming method for displaying a map of a plane (ground) on the basis of a bird\"\"s-eye view corresponding to a projected plan of the plane (ground) which is obtained by projecting the plane (ground) on any plane between the plane (ground) and a view point when the plane (ground) is viewed from the view point located at any height from the plane (ground), a map display apparatus using the bird\"\"s-eye view forming method, and a navigation system using the map display apparatus.\nIn order to attain the above object, a bird\"\"s-eye view forming method for forming, from map data, drawing (display) data for a map which is to be represented by a bird\"\"s-eye view, is characterized in that coordinate data contained in the map data are perspectively converted to drawing (display) data based on a bird\"\"s-eye view on a desired projection plane with a view point being set to a desired position (i.e., the coordinate data of the map data on a vertical orthogonal projection plane are perspectively converted to drawing (display) data on a desired projection plane based on a bird\"\"s-eye view).\nThe bird\"\"s-eye view forming method preferably includes a step of receiving an input of the position of the view point, and a step of determining the projection plane so that the drawing (display) positions of two predetermined spots (for example, a current position and a destination) which are obtained through a perspectively projecting transformation on the basis of the coordinates of the two spots in the map data and the position of the view point, are set to predetermined positions. (Here, the perspectively projecting conversion is defined as such a data conversion that data on a plan view are converted to data on a bird\"\"s-eye view through a perspective projection operation. Therefore, by using the perspectively projecting conversion, a plan-view map is converted to a bird\"\"s-eye view map. In the following description, the perspectively projecting conversion is merely referred to as xe2x80x9cperspective projectionxe2x80x9d).\nAlternatively, the bird\"\"s-eye view forming method may include a step of receiving an input of a scale, and a step of determining the position of the view point and the projection plane so that the drawing (display) positions of the two predetermined spots (for example, a current position and a destination) which are obtained by performing the perspective projection on the basis of the coordinates of the two spots in the map data and the input scale, are set to predetermined positions, and a scale of a drawn (displayed) map is equal to the input scale.\nFurther, the bird\"\"s-eye view forming method may include a step of receiving an input of a projection angle which is defined as an angle at which a plane defined for the map data and the projection plane intersect each other, and a step of determining the projection plane on the basis of the input projection angle and the set view point position.\nStill further, according to the present invention, there are provided a map display apparatus for displaying a bird\"\"s-eye view by using the bird\"\"s-eye view forming method as described above, and a navigation system for displaying a map by using the map display apparatus.\nAccording to the bird\"\"s-eye view display apparatus and the navigation system, map drawing means includes coordinate transforming means, and the coordinate transforming means performs the perspective projection (conversion) on the map data to convert a plan-view map to a bird\"\"s-eye view map, and to display the converted map on a screen. Therefore, according to the present invention, a user can obtain a bird\"\"s-eye view display which is easy to see and in which the user can easily recognize the positional relationship of spots displayed on the map. According to the bird\"\"s-eye view forming method of the present invention, the view point can be freely set, so that it can meet user\"\"s need.\nAccording to the present invention, a navigation system which is excellent in operational performance and convenient for a user can be obtained if it is designed to overcome the following first to seventh problems.\nFirstly, in the case where the processing speed of scroll processing in a bird\"\"s-eye view display mode is lower than that in a plan display mode, the operational performance may be degraded when, the scroll processing is frequently used to search for a desired spot on the bird\"\"s-eye view display (first problem).\nTherefore, it is preferable in the present invention that the plan view (plan map) display mode and the bird\"\"s-eye view display mode are freely switchable. With this switching operation, both the plan view display and the bird\"\"s-eye view display can be made in accordance with a user\"\"s requirement, and thus the convenience is enhanced.\nIn the above case, a second problem occurs, namely, that it is difficult for the user to grasp the positional relationship between two different type of view when display is switched between the plan view display and the bird\"\"s-eye view display. This is because the position of the same spot to be displayed on the screen varies significantly due to variation of the view point with respect to the map, when the display switching operation is performed, and further, spots which have not been displayed or have been displayed until the display switching time are respectively displayed or are not displayed at the display switching time.\nTherefore, if the present invention is designed so that the display is freely switchable the plan view display and the bird\"\"s-eye view display, it is further preferable that an intersecting angle between a plane containing the map data and a plane onto which the map is projected, is gradually increased or reduced in time series. That is, according to the present invention, it is preferable that the view point is smoothly shifted during the conversion between the plan view display and the bird\"\"s-eye view display. With this operation, the shift between the plan view and the bird\"\"s-eye view is smoothly performed, and thus the user can easily recognize the positional relationship between spots displayed on the plan view map and the same spots displayed on the bird\"\"s-eye view map.\nThirdly, when the map is displayed in the bird\"\"s-eye view display mode, if the view point serving as a parameter to the perspective projection is fixed to a position, a mark indicating the current position which is displayed on the screen is shifted on the map while the current position is shifted. In this case, if the marked current position deviates from a displayed map area, the mark of the current position is not displayed (third problem).\nTherefore, according to the present invention, the view point from which to obtain the bird\"\"s-eye view may be fixed to a specific position which is set to be away from the current position at a fixed distance and in a fixed direction at all times. Furthermore, the height of the view point may be varied in accordance with an external operation. If the view point is fixed to the specific position as described above, the bird\"\"s-eye view is displayed such that the current position is fixed to a point on the screen at all times, so that the user can easily grasp the current position. Furthermore, if the position of the view point can be set to a user\"\"s desired position, the convenience can be further enhanced.\nFourthly, when two spots are displayed on the same frame, it is sufficient in the conventional plan view display to renew the scale in accordance with the distance between the two points. However, in the bird\"\"s-eye view display of the map, the map cannot be optimally displayed by merely renewing the scale (fourth problem).\nTherefore, according to the present invention, an instruction regarding the positions of two spots (for example, the current position and the destination) may be accepted, and the view point and the intersection angle between the map-projected plane (bird\"\"s-eye view) and the plane containing the map data (plan view) may be determined so that the two spots are displayed at predetermined positions on the screen, whereby the bird\"\"s-eye view display is performed. With this operation, the positional relationship between the two spots can be easily recognized. Further, even when the positions of the two spots vary, these are displayed on the same frame at all times, so that the user can grasp the positional relationship between the two spots with no complicated operation.\nAccording to the present invention, in the bird\"\"s-eye view display mode, line backgrounds such as roads, railroad lines, etc., plane backgrounds such as rivers, green zones, etc., and character strings may be perspectively converted and drawn on the bird\"\"s-eye view map. However, if the character strings are perspectively converted, the shape of a character is smaller and is more deformed as it becomes farther away from the view point. On the other hand, if it is in the vicinity of the view point, it is enlarged. Thus the character strings may be illegible in these cases (fifth problem).\nTherefore, it is preferable that the coordinate transforming means of the present invention does not perform the perspective projection on a character image. In this case, characters contained in the bird\"\"s-eye view are displayed at the same size, so that the character strings are easily legible.\nFurther, when a lot of character data are displayed, character strings are displayed while overlapping with one another. Particularly in the bird\"\"s-eye view, an angle of depression with respect to the view point is smaller at a far-distance point from the view point and thus the reduction scale is large at that point, so that the data amount of character data to be displayed per unit area is increased. Accordingly, the character data (strings) are liable to be overlapped with each other at the far-distance point (sixth problem).\nTherefore, it is preferable that as the position of character data contained in the map data becomes nearer to the view point, drawing (display) judgment means of the present invention displays the character data more preferentially (i.e., increases the display priority rank of the character data). With this operation, character data having a higher display priority rank may be displayed while being superposed on character data having a lower display priority rank, so that the character string can be prevented from being deficient and it can be made legible.\nWhen the current position is located in the vicinity of the bottom side of the display frame or the view point, it is preferable that the display height or the distance from the view point is set as a criterion for determining the display priority rank. Further, when two or more character strings are displayed while overlapping with one another, a character string which is displayed at a lower display height (which is the height from the bottom side of the display frame) on the screen may be displayed more preferentially (i.e., so that the character string is not covered by the other character strings, but is superposed on the other character strings). In this case, the character data near to the current position are displayed while being prevented from being covered by the other character strings and thus missed. With this display operation, the deficiency (missing) of the character information (string) near to the current position which the user has a greater need for can be prevented, and the character string can be made legible.\nAs described above, when the bird\"\"s-eye view display is performed, the data amount of line data, plane background data and character data to be drawn per unit area is increased in a region which is far away from the view point (hereinafter referred to as xe2x80x9cfar-distance regionxe2x80x9d, and this region is a region having a small angle of depression). Accordingly, when the character data of these data are drawn while superposed on the other data, the line data and the plane background data are covered by the character data in the far-distance region, so that it may be difficult to check the shapes of roads, etc. in the far-distance region (seventh problem).\nTherefore, the drawing judgment means of the present invention preferably eliminates, from targets to be drawn (displayed), those character strings which are obtained through the perspective projection of the map and which are located at positions higher than a predetermined display height from the bottom side of the display frame. Or the drawing judgement means may eliminate those character strings which are defined at position where is further than a fixed distance from the view point. As described above, the character strings are not displayed (drawn) in a region which is located above the predetermined height or is more than at a predetermined distance away, that is, in a region where the data amount to be drawn per unit area is considerably increased due to the reduction in angle of depression in the bird\"\"s-eye view display mode. Therefore, a road display, a plane background display, etc. are not covered by the character string display, so that the recognition of roads in the far-distance regions can be prevented from being impaired even in the bird\"\"s-eye view display mode."}
{"text": "This invention relates to a process for the treatment of strongly acidic cation exchange resin catalysts based on styrene/divinyl benzene copolymerizates.\nRecently, cation exchange resin catalysts have been employed in processes for the ecologically beneficial implementation of acid-catalyzed syntheses. The syntheses in question are for instance esterifications, splitting of esters, hydrolyses, condensations, hydrations as well as alkylations and acetylations of aromatics. Unlike liquid acids, they have the advantage that the catalyst can be easily separated from the product and that no waste acid mixtures are produced as in the conventional homogeneous catalysis.\nA prerequisite for the practical use of a solid cation exchanger instead of a liquid acid is, in addition to sufficient selectivity and space/time yield, the thermal stability of the copolymerizates under the respective reaction conditions.\nStrongly acidic styrene/divinyl benzene copolymerizates, the core of which has been substituted with halogen, are particularly thermally stable and are used at temperatures in the range of 100.degree. to 200.degree. C. for acid-catalyzed syntheses, such as for the hydration of lower olefins or for alkylation reactions."}
{"text": "In many camera surveillance situations it is important to detect activity, such as presence of objects or motion, at large distances. For natural reasons, objects which are located far away from the camera will be small in images captured by the camera. For instance, they may only cover a few pixels in the images. In order to detect activity at large distances, it is thus necessary to, e.g., be able to detect objects which appear as small in the captured images.\nHowever, detecting small objects usually comes at the cost that false detections increase. This has to do with the fact that there is noise in the images captured by the camera. The noise may easily be mistaken for being a small object and vice versa, thereby leading to false detections. In order to reduce the false detections, video processing algorithms for object detection, motion detection, and/or object tracking, include parameters for such things as intensity change thresholds and noise filtering, including morphological filters to remove small detections. The parameters are typically set so that false detections caused by noise are removed while there is still enough sensitivity to detect actual activity. In other words, there is a trade-off between detection range and avoiding false detections.\nThe problem of detecting and identifying objects at a large distance has been recognized. In order to solve the problem, it has been suggested to use a thermal camera to identify a position coordinate of a person. The position coordinate is then input to a control method for controlling a second camera, and in particular the direction and zoom thereof, to capture an enlarged image of the person.\nSuch a method of has the drawback that the second camera is controlled to change its direction and zoom level to follow a specific target object. Accordingly, such method is only suitable for detecting one object at a time.\nThere is thus room for improvements. In particular, there is a need for alternative methods which allow detection at large ranges at the same time as false detections are reduced."}
{"text": "A considerable number of solar heat exchange panels presently exist. In a solar heating system, the solar heat exchange panel, which is the heat collecting portion, is exposed to the sun to absorb heat from the rays of the sun. Air is circulated within the heat exchange panel to transfer the heat from the solar heat exchange panel to where the heat can be utilized. Frequently a liquid conduit is also mounted in the solar heat exchange panel to absorb heat which can also be used for various purposes as is the heated air.\nThis invention relates to improvements in the structure of a solar heat exchange panel and in particular to a solar heat exchange panel of limited size to relationship to the amount of heat that it can absorb. Prior art relating to solar heat exchange panels is as follows:\n______________________________________ Inventor Patent Number ______________________________________ Moreau 201,439 Awot 2,907,318 Gallagher 4,011,856 Taylor 4,016,861 ______________________________________"}
{"text": "The use of chromium compounds in the polymerization of olefins is well-known. U.S. Pat. Nos. 2,825,721 and 2,951,816 teach the use of CrO.sub.3 supported on an inorganic material such as silica, alumina or combinations of silica and alumina and activated by heating at elevated temperatures to polymerize olefins. When these catalyst systems are used in various polymerization processes such as the well-known particle-form process, the resins produced, while useful in many applications, are unsatisfactory for others because of a deficiency in certain properties such as melt index.\nAttempts to improve deficient properties of polyolefins produced using supported, heat-activated chromium oxide catalysts have been made by adding various metal and non-metal compounds to the supported chromium oxide prior to activation by heating. For example, in U.S. Pat. No. 3,622,522 it is suggested that an alkoxide of gallium or tin be added to supported chromium oxide prior to heat activation, and the performance of aluminum isopropoxide compared. U.S. Pat. No. 3,715,321 suggests adding a compound of a Group II-A or Group III-B metal to supported chromium oxide prior to heat treatment.\nU.S. Pat. No. 3,484,428 discloses an olefin polymerization catalyst composite prepared by impregnating a support such as silica, alumina, zirconia, and thoria with a chromium compound, convertible to chromium oxide upon calcination, and a borane or an alkyl-substituted borane, and then heating and activating both the boron and chromium compounds of the resulting composite on the support in an oxygen-containing atmosphere at a temperature of about 750.degree. to 1800.degree. F. U.S. Pat. No. 3,780,011 discloses an olefin polymerization catalyst prepared by forming an intimate mixture of a silica, alumina, zirconia or thoria support and chromium oxide. The mixture is dried by fluidizing with dry air. An oxygenated titanium, boron or vanadium compound, e.g., alkyl esters of such metals is then added to the fluidized bed, and the resulting fluidized mixture is activated with an oxygen-containing, dry gas at elevated temperatures.\nU.S. Pat. No. 3,202,645 defines catalysts, including the combination of a halide of titanium, zirconium, chromium, etc. with the reaction product of e.g., tri-n-butyl borane with a particulate inorganic solid such as silica. U.S. Pat. No. 3,625,864 relates to a supported catalyst including a titaniaboron complex. U.S. Pat. No. 2,898,326 describes a catalyst comprising an activated Group VI-A oxide support impregnated with boranes.\nIt is also known to utilize other chromium compounds as catalysts for the polymerization of olefins. Such compounds include various silyl chromate and polyalicyclic chromate esters as described, for example, in U.S. Pat. Nos. 3,324,095; 3,324,101; 3,642,749; and 3,704,287. The use of phosphorus-containing chromate esters in olefin polymerization catalysts has also been disclosed in the aforesaid U.S. Pat. No. 3,704,287 and in U.S. Pat. No. 3,474,080.\nUse of the above chromium compound catalysts in Ziegler-type coordination catalyst systems has also been proposed. As is well-known in the art, such catalysts frequently additionally comprise organometallic reducing agents such as, for example, trialkyl aluminum compounds. Ziegler-type catalyst systems incorporating supported chromium compound catalysts and organometallic reducing agents, particularly organoaluminum compounds, are disclosed, for example, in U.S. Pat. Nos. 3,324,101; 3,642,749; 3,704,287; 3,806,500; and in the copending application Ser. No. 532,131, filed Dec. 16, 1974.\nThe preparation and use of improved high pore volume silica xerogel materials suitable as catalyst supports is described in Belgium Pat. No. 741,437 and U.S. Pat. Nos. 3,652,214; 3,652,215 and 3,652,216; whereas moderate pore volume silica xerogel supports are disclosed in U.S. Pat. No. 3,453,070."}
{"text": "Computers normally include an arithmetic logic unit (ALU) that includes an adder that adds numbers of some maximum number of bits. Adders for words of length 32 and 64 bits are common in microprocessors and the like. While these adders will also operate on much smaller words, when doing so, the majority of the logic circuits contained in the adders are idle. For example, a 64-bit adder can be used to add two 8-bit words by placing each of the 8-bit words in the least significant portion of a corresponding 64 -bit word and then adding the 64 -bit words. During the addition, the logic circuitry concerned with adding the 7 high order bytes of each of the words is effectively idle. Hence, 7/8.sup.ths of the capacity of the adder is being wasted during this operation.\nComputations involving operations on a large number of small words are often encountered in multi-media data processing. Images are normally represented as arrays of pixels in which each pixel is represented by a word that is significantly smaller than maximum word size of the arithmetic logic unit. A gray scale image is typically represented by an array of one byte integers representing the light intensity at corresponding locations in the image. Similarly soundtracks are typically represented by arrays of one or two byte integers representing the intensity of the sound track as a function of time. Hence, multi-media data processing typically under utilizes the computational capacity of the arithmetic logic unit incorporated in the typical general purpose computer.\nIn addition to under utilizing the capacity of the ALU, this type of data processing often results in further inefficiencies resulting from the need to pack and unpack the data prior to operating on it in the ALU. For example, since storage space is always at a premium, the pixels of the image are typically packed into words. If the basic word size on the computer is 32-bits, the pixels of a gray scale image could be packed four per word. Consider an operation that must be performed on each pixel in the image. In addition to the time needed to perform the operation, the program must also unpack the pixel information prior to each computation and repack the result. These packing and unpacking operations further decrease the efficiency of the operation.\nThe computation times encountered in multi-media operations can be excessive. Hence, special parallel computer architectures are often employed to reduce the time between the execution of the sum image command and the time at which the sum image is completed. Multimedia processing operations performed on one pixel or sound sample are often independent of the operations being performed on the other pixels or sound samples; hence, the operations can be performed in parallel without regard to ordering. A computer with M adders can, in principle, provide a result in 1/M.sup.th the time provided the movement of the pixels between memory and the adders does not become a bottleneck. Hence, it would be advantageous to provide a computer architecture in which multiple additions can be performed in parallel. Unfortunately, the cost of providing these additional ALUs and the hardware needed to control them is often prohibitive.\nOne class of computation that is often performed over a large data set is multiplication by a binary fraction. Such computation are common in filtering operations and in data compression and decompression operations. In these cases, the multiplication is by a constant.\nBroadly, it is the object of the present invention to provide an improved ALU.\nIt is a further object of the present invention to provide an ALU that operates at high efficiency when multiple operations involving words that are smaller than the width of the ALU are processed.\nIt is a still further object of the present invention to provide an ALU that is adapted for computing a plurality of multiplications involving a binary fraction.\nThese and other objects of the present invention will become apparent to those skilled in the art from the following detailed description of the invention and the accompanying drawings."}
{"text": "1. Field of the Invention\nThis invention relates to fixtures of the bed of nails type used to electrically connect a printed circuit board to be tested with an automatic test system.\n2. Description of the Prior Art\nAs is well known, it is becoming a routine matter among manufacturers to test circuit boards which they intend to use in producing products such as televisions, electronic organs, and many, many others. These tests can be conducted with so-called \"bare\" boards and with \"loaded\" boards. A bare board is one not having components mounted thereon while a loaded board is one having these components in place.\nIn conducting such tests, it is, of course, necessary to have a device commonly called a \"fixture\" which serves to make an electrical connection between the nodal points in the circuit to be tested and the test system. A standard approach is to use a \"bed of nails\" fixture. This device derives its name from the fact that it is comprised of a great number of nail-like test probes having tips to make electrical contact with the nodal points of the circuit to be tested. These probes are individually connected through one or more interfacing devices to the test system.\nIt has been the practice for manufacturers to purchase one fixture for each type of board to be tested. This has been primarily due to the fact that fixtures have largely been custom-made. One would lay out the pattern of nodal points to be used in a test and an array of test probes would be provided in that pattern in one fixture. As the number of boards one must test grows, this becomes an expensive proposition. The expense is due primarily to the high number of probes involved. Test probes are typically spring-loaded and gold-plated and are manufactured to fairly close tolerances. It is not difficult to imagine, then, that a substantial investment is involved if a manufacturer, as is now the practice, is required to maintain an inventory of several fixtures, each containing many such probes.\nIt would, therefore, be advantageous to provide a universal fixture which could be used to test many different types of circuit boards. It can be expected that such a fixture would represent a relatively large initial investment because it would need to contain more probes than would ordinarily be required to test a single board. In the long run, however, this approach would prove to be a more economical one for manufacturers who are required to test many types of boards."}
{"text": "Embodiments of the present invention relate to a semiconductor device and a method of manufacturing the same and, more particularly, to a three-dimensionally (3-D) structured nonvolatile memory device and a method of manufacturing the same.\nA nonvolatile memory device retains stored data although the supply of power is blocked. As 2-D structured memory devices including memory cells fabricated in a single layer on a silicon substrate reach the limit in increasing the integration degree thereof, there is proposed a 3-D structured nonvolatile memory device in which memory cells are vertically stacked on a silicon substrate.\nThe structure and features of a conventional 3-D nonvolatile memory device are described below with reference to FIG. 1.\nFIG. 1 is a cross-sectional view illustrating the structure of a conventional 3-D nonvolatile memory device.\nAs shown in FIG. 1, the conventional 3-D nonvolatile memory device includes channels CH protruding from a substrate 10 and a plurality of memory cells MC vertically stacked along the channels CH. The memory device further includes a lower selection gate LSG formed under the plurality of memory cells MC and an upper selection gate USG formed over the plurality of memory cells MC. Bit lines BL are provided over the upper selection gate USG and coupled to the channels CH. In this structure, a plurality of memory cells MC coupled in series between the lower selection gate LSG and the upper selection gate USG form one a cell string STRING, and the cell strings STRING are arranged on the substrate 10.\nIn FIG. 1, reference numerals 11, 14, and 17 denote interlayer insulating layers, reference numeral 12 denotes a lower selection line, reference numerals 15 denote word lines, and reference numeral 18 denotes an upper selection line. Furthermore, reference numerals 13 and 19 denote gate insulating layers, reference numeral 16 denotes a charge blocking layer, a charge trap layer, and a tunnel insulating layer.\nA method of forming the memory cells MC is described below in short. First, after alternately forming a plurality of conductive layers and a plurality of interlayer insulating layers, trenches are formed by etching the plurality of conductive layers and the plurality of interlayer insulating layers. After forming the charge blocking layer, the charge trap layer, and the tunnel insulating layer 16 on the inner walls of the trenches, a channel layer is filled within the trenches. In this manufacture process, the charge trap layers of the plurality of memory cells MC stacked along each of the channels CH are coupled.\nHere, the charge trap layer functions as a substantial data depository for storing data through the injection or discharge of electric charges into or from the data depository. Accordingly, in the conventional structure in which the charge trap layers of the memory cells MC are coupled, data stored in memory cells MC may be lost because electric charges stored in one memory cell MC move to another memory cell MC. For example, if the charge trap layer is formed of a Si-rich nitride layer, stored data may be lost because electric charges stored in the charge trap layer move. In order to prevent the stored data from being lost, the charge trap layer may be made of stoichiometric nitride. If the charge trap layer is made of stoichiometric nitride, however, an erase operation speed may be slow."}
{"text": "1. Field of the Invention\nThe present invention relates to a system, method, and program for selecting a path to use to communicate to a device to improve transmission performance.\n2. Description of the Related Art\nTwo systems communicating over a network may each include multiple ports, thus providing multiple paths across which data can be communicated. In certain prior art systems, a path may be selected according to a round robin or other predefined path rotation technique or a single default path is used for all operations. However, such techniques do not attempt to optimize path selection when there are multiple available paths.\nFor this reason, there is a need in the art to provide a methodology for selecting paths that improves data transfer rates across the system when there are multiple paths between the two systems."}
{"text": "Modern lubricating oil formulations are formulated to exacting specifications often set by original equipment manufacturers. To meet such specifications, various additives are used, together with base oil of lubricating viscosity. Depending on the application, a typical lubricating oil composition may contain dispersants, detergents, anti-oxidants, wear inhibitors, rust inhibitors, corrosion inhibitors, foam inhibitors just to name a few. Different applications will govern the type of additives that will go into a lubricating oil composition.\nA functional fluid is a term which encompasses a variety of fluids including but not limited to tractor hydraulic fluids, automatic transmission fluids including continuously variable transmission fluids, manual transmission fluids, hydraulic fluids, power steering fluids, fluids related to power train components and fluids which have the ability to act in various different capacities. It should be noted that within each of these fluids such as, for example, automatic transmission fluids, there are a variety of different types of fluids due to the various transmissions having different designs which have led to the need for fluids of markedly different functional characteristics.\nTractor hydraulic fluids and automatic transmission fluids are examples of functional fluids having very specific friction requirements. Because such fluids work in wet brake and/or wet clutch systems, the fluid must assist in smooth engagement of these brakes and clutches while maintaining desirably high frictional properties for effective brakes and clutches. These fluids require high friction coefficients. For example, tractor hydraulic fluids that involve wet brake systems must have a high friction coefficient to be effective. Further, automatic transmission fluids must have enough friction for the clutch plates to transfer power. However, the friction coefficient of fluids has a tendency to decline due to the temperature effects as the fluid heats up during operation. It is important that the tractor hydraulic fluid or automatic transmission fluid maintain its high friction coefficient at elevated temperatures, otherwise brake systems or automatic transmissions may fail."}
{"text": "The polypeptide insulin is the primary hormone responsible for controlling the transport, utilization and storage of glucose in the body. The β-cells of the pancreatic islets secrete a single chain precursor of insulin, known as proinsulin. Proteolysis of proinsulin results in removal of certain basic amino acids in the proinsulin chain and the connecting or C-peptide and provides the biologically active polypeptide insulin.\nThe insulin molecule has been highly conserved in evolution and generally consists of two chains of amino acids linked by disulfide bonds. In the natural human, two-chain insulin molecule (mw 5,800 Daltons), the A-chain is composed of 21 amino acid residues and has glycine at the amino terminus; and the B-chain has 30 amino acid residues and phenylalanine at the amino terminus.\nInsulin may exist as a monomer or may aggregate into a dimer or a hexamer formed from three of the dimers. Biological activity, i.e., the ability to bind to receptors and stimulate the biological actions of insulin, resides in the monomer.\nDiabetes is a biological disorder involving improper carbohydrate metabolism. Diabetes results from insufficient production of or reduced sensitivity to insulin. In persons with diabetes, the normal ability to use glucose is inhibited, thereby increasing blood sugar levels (hyperglycemia). As glucose accumulates in the blood, excess levels of sugar are excreted in the urine (glycosuria). Other symptoms of diabetes include increased urinary volume and frequency, thirst, itching, hunger, weight loss, and weakness.\nThere are two varieties of diabetes. Type I is insulin-dependent diabetes mellitus, or IDDM. IDDM was formerly referred to as “juvenile onset diabetes.” In IDDM, insulin is not secreted by the pancreas and must be provided from an external source. Type II or adult-onset diabetes can ordinarily be controlled by diet, although in some advanced cases insulin is required.\nBefore the isolation of insulin in the 1920s, most patients died within a short time after onset. Untreated diabetes leads to ketosis, the accumulation of ketones, products of fat breakdown, in the blood. This is followed by the accumulation of acid in the blood (acidosis) with nausea and vomiting. As the toxic products of disordered carbohydrate and fat metabolism continue to build up, the patient goes into a diabetic coma, which leads to death.\nThe use of insulin as a treatment for diabetes dates to 1922, when Banting et al. (“Pancreatic Extracts in the Treatment of Diabetes Mellitus,” Can. Med. Assoc. J., 12:141-146 (1922)) showed that the active extract from the pancreas had therapeutic effects in diabetic dogs. In that same year, treatment of a diabetic patient with pancreatic extracts resulted in a dramatic, life-saving clinical improvement.\nUntil recently, bovine and porcine insulin were used almost exclusively to treat diabetes in humans. Today, however, numerous variations in insulin between species are known. Each variation differs from natural human insulin in having amino acid substitution(s) at one or more positions in the A- and/or B-chain. Despite these differences, most mammalian insulin has comparable biological activity. The advent of recombinant technology allows commercial scale manufacture of human insulin (e.g., Humulin™ insulin, commercially available from Eli Lilly and Company, Indianapolis, Ind.) or genetically engineered insulin having biological activity comparable to natural human insulin.\nTreatment of diabetes typically requires regular injections of insulin. Due to the inconvenience of insulin injections, massive efforts to improve insulin administration and bioassimilation have been undertaken.\nAttempts have been made to deliver insulin by oral administration. The problems associated with oral administration of insulin to achieve euglycemia in diabetic patients are well documented in pharmaceutical and medical literature. Digestive enzymes in the gastrointestinal tract rapidly degrade insulin, resulting in biologically inactive breakdown products. In the stomach, for example, orally administered insulin undergoes enzymatic proteolysis and acidic degradation. Comparable proteolytic breakdown of insulin occurs in the intestine. In the lumen, insulin is attacked by a variety of enzymes including gastric and pancreatic enzymes, exo- and endopeptidases, and brush border peptidases. Even if insulin survives this enzymatic attack, the biological barriers that must be traversed before insulin can reach its receptors in vivo may limit its bioavailability after oral administration of insulin. For example, insulin may possess low membrane permeability, limiting its ability to pass from the intestinal lumen into the bloodstream.\nSome efforts to provide an oral form of insulin have focused on providing insulin-oligomer conjugates. Human insulin and many closely related insulins that are used therapeutically contain three amino acid residues bearing free primary amino groups. All three primary amino groups, namely the N-termini (alpha amino groups) of the A and B chains (GlyA1 and PheB1) and the epsilon-amino group of LysB29, may be modified by conjugation with oligomers. Depending on the reaction conditions, N-acylation of an unprotected insulin leads to a complex mixture of mono-, di-, and tri-conjugates (e.g., insulin mono-conjugated at GlyA1, insulin mono-conjugated at PheB1, insulin mono-conjugated at LysB29, insulin di-conjugated at GlyA1 and PheB1, insulin di-conjugated at GlyA1 and LysB29, insulin di-conjugated at PheB1 and LysB29, and insulin tri-conjugated at GlyA1, PheB1, and LysB29). When a particular conjugate, for example insulin mono-conjugated at LysB29, is desired, it can be burdensome and/or expensive to separate (or purify) such a complex mixture of conjugates to obtain the desired conjugate.\nAs a result, various efforts have been undertaken to selectively synthesize the desired insulin conjugate. For example, Muranishi and Kiso, in Japanese Patent Application 1-254,699, propose a five-step synthesis for preparing fatty acid insulin derivatives. The A1- and B1-amino groups of insulin are protected (or blocked) with p-methoxybenzoxy carbonyl azide (pMZ). After acylation with a fatty acid ester, the protection (blocking) groups are removed to provide insulin mono-acylated at Lys(B29) with a fatty acid. As another example, U.S. Pat. No. 5,750,497 to Havelund et al. proposes treating human insulin with a Boc-reagent (e.g. di-tert-butyl dicarbonate) to form (A1,B1)-diBoc human insulin, i.e., human insulin in which the N-terminal end of both the A- and B-chains are protected by a Boc-group. After an optional purification, e.g. by HPLC, a lipophilic acyl group is introduced in the ε-amino group of LYSB29 by allowing the product to react with a N-hydroxysuccinimide ester of the formula X—OSu wherein X is the lipophilic acyl group to be introduced. In the final step, trifluoroacetic acid is used to remove the Boc-groups and the product, NεB29—X human insulin, is isolated.\nVarious other efforts have been undertaken to preferentially synthesize the desired insulin conjugate to provide a mixture of conjugates in which the desired insulin conjugate is the preferred product. For example, U.S. Pat. No. 5,646,242 to Baker et al. proposes a reaction that is performed without the use of amino-protecting groups. Baker proposes the reaction of an activated fatty ester with the ε-amino group of insulin under basic conditions in a polar solvent. The acylation of the ε-amino group is dependent on the basicity of the reaction. At a pH greater than 9.0, the reaction preferentially acylates the ε-amino group of B29-lysine over the α-amino groups. Examples 1 through 4 report reaction yields of the mono-conjugated insulin as a percentage of the initial amount of insulin between 67.1% and 75.5%. In Example 5, Baker also proposes acylation of human proinsulin with N-succinimidyl palmitate. The exact ratios of ε-amino acylated species to α-amino acylated species were not calculated. The sum of all ε-amino acylated species within the chromatogram accounted for 87-90% of the total area, while the sum of all related substances (which would presumably include any α-amino acylated species) accounted for <7% of the total area, for any given point in time.\nIt is desirable to provide methods of site specifically synthesizing desired, particular insulin-oligomer conjugates that may be less burdensome and/or more cost effective than the conventional methods described above."}
{"text": "1. Field of the Invention\nThe invention relates to a system for monitoring a plurality of capacitor cells connected in series in a capacitor module which are used as electrical energy stores in an energy storage device and which can be charged in same by a charge current flowing through all the capacitor cells or discharged by a discharge current.\n2. Background Information\nCapacitor cells connected in series of this type have the problem that their capacity can change due to aging processes or that parasitic discharging can occur and thus disproportionate voltages can ensue for the individual capacitor cells over time which ultimately have a negative impact on the capacitor module's capacity to store energy.\nThe object of the invention is thus to provide a system for monitoring such capacitor cells which prevents such problems caused by changes to the capacity of individual cells."}
{"text": "On Mar. 29, 1984, the FCC Report and Order in Docket No. 21323 adopted and authorized a standard for multichannel television sound (MTS). Bulletin No. 60 of the Office of Science and Technology, \"Multichannel Television Sound Transmission and Audio Processing Requirement for the BTSC System\" (OST 60) contains the technical specifications for the above mentioned MTS adopted by the FCC.\nThe MTS system, known variously as MTS, BTSC stereo, or broadcast TV stereo, consists of two parts: an overall transmission system, and a noise reduction, or companding system. The \"BTSC System Multichannel Television Sound Recommended Practices\" (BTSC Standard) defines companding as: \"a noise reduction process used in the stereophonic subchannel and in the second audio program subchannel consisting of compression (encoding) before transmission and complementary expansion (decoding) after reception. (This definition conforms to \"companding\" in OST 60 Section A). \"\nWithout companding, the transmission system is capable of delivering high-quality stereophonic audio. However, FM-transmission systems experience squared relationship between noise and frequency (Parabolic Noise Characteristic) resulting in higher noise at the higher frequencies. In addition, to avoid causing interference, the BTSC Standard limits the amount of modulation that can be applied to the stereophonic signal. Thus, even under ideal conditions, the addition of the subcarrier adds approximately 15 dB of noise to stereophonic (stereo) reception compared to monophonic (mono) reception. To make matters worse, under certain impaired transmission/reception conditions, such as a weak received signal, transmitter ICPM and multipath effects, buzz or hum can be introduced onto the transmitted audio. Thus, without companding, the service area for stereophonic TV (stereo-TV) reception is smaller than for monophonic TV (mono-TV) reception.\nThe situation is worse for the Second Audio Program (SAP) channel. The SAP subcarrier frequency is 78.7 KHz, much higher than for the stereo signal; therefore, the Parabolic Noise Characteristic results in even more noise in the signal. Furthermore, this subcarrier uses FM, which is additionally subject to picture-to-audio intermodulation (buzz beat), causing a particularly annoying distortion.\nThe BTSC noise-reduction system was designed to be a cost-effective aid to the MTS transmission system in delivering a clean, noise-free audio signal into the home. Specifically, the system was designed to: provide significant noise reduction even in poor reception areas while preserving input-signal dynamic range; prevent the stereo-subcarrier from interfering with overall transmitted power levels (AM-interleave effects); ensure reliable performance even in the face of manmade noise and transmission/reception-system impairments; and to provide these benefits at a commercially reasonable cost.\nTo achieve the above-stated goals, the BTSC Standard departed from previous design approaches where the dynamic range of the impaired channel was very low. In the case of the stereo-subcarrier in Grade B reception, the available dynamic range was about 43 dB, while in the SAP channel the dynamic range was about 26 dB. These compare unfavorably with the typical dynamic range of a compact cassette at about 60 dB. The significance of these figures, the operation of prior art compander systems, and the operation of the current invention will be better understood after a discussion of the psychoacoustic phenomenon of masking.\nAll audio noise-reduction systems work on the principle of masking; a listener will be oblivious to the noise on a transmission when the program signal, music or speech, is loud enough and its spectral content is broad enough, to mask the noise.\nFor example, if the program consists of low-frequency sounds, it must be transmitted at a high level relative to the background noise of the stereo-subcarrier channel to capture the listener's attention and for the listener to be unaware of the background noise. On the other hand, a broadband signal does not need to be much higher in amplitude than the background noise for the noise to fade below the listener's perception threshold. See, for example, I. M. Young, and C. H. Wenner, Masking of White Noise by Pure Tone, Frequency-Modulated Tone, and Narrow-Band Noise, J. Acoust. Soc. Am. 41, pp. 700-705, 1967.\nThe noise reduction system must compress and encode the audio signal such that it will consistently mask the channel noise during transmission, and then decode and expand the transmitted signal to recover the original audio signal. In passing through the encode/decode (companding) cycle, distortion or other degradation of the audio signal must be kept to a minimum. And in the decoding process, all the audible noise should be eliminated. Thus, not only must the level of the transmitted audio be high relative to the background noise, but the frequency of the signal and noise must be considered when selecting or designing an effective noise masking and companding scheme. Other characteristics of the signal and noise will affect the design of an ideal companding system. Thus, the amplitude of the signal, the rate at which the signal changes amplitude, and even whether the signal is decreasing or increasing are also important parameters in the design of an ideal companding system.\nThe stereo-subcarrier's background-noise spectrum is white, rising at 3 dB-per-octave. By comparison, the SAP subcarrier noise rises at 9 dB-per-octave. The noise can be masked if the transmitted signal's spectrum contains substantial high-frequency, especially in the case of the SAP channel. If so, the compander would only have to keep the signal amplitude levels high through the transmission medium. However, most TV program materials have their dominant energy at low frequencies. Alternatively, if the program consistently lacked high-frequency content, one could simply apply a constant rising preemphasis characteristic to the entire audio spectrum.\nHowever, today's TV program content, especially music and movie effects, is neither consistently high frequency nor low frequency, for either approach to work. Inevitably, the signal will have instances of high-level, high-frequency energy in the audio signal, where a constant rising preemphasis would cause headroom (overload distortion) problems.\nExisting solutions to the problem of preserving psychoacoustic masking and simultaneously preserving headroom at all frequencies use spectral companding, a preemphasis scheme which adapts its characteristic to suit the signal. The spectral compressor in the encoder measures the spectral balance of the input signal and varies the high-frequency-preemphasis accordingly, merely increasing the potential for masking, and reducing the possibility of high-frequency overload. The resulting encoded signal is, therefore, dynamically adjusted to consistently contain a substantial proportion of high frequencies before transmission, thereby masking the channel noise.\nDuring reception, the spectral expander (in the decoder) restores the high frequencies to their proper amplitude. If the original input signal contains predominantly low frequencies, the decoder attenuates the high-frequency background noise, leaving the low-frequency signal and low-frequency background noise, the latter of which is masked by the signal itself. If the original input signal contains predominantly high frequencies, the decoder does not attenuate the high frequencies to restore correct frequency response, since the signal itself masks the noise.\nThe spectral compressor achieves two simultaneous requirements: the system is forgiving of high-background-noise environments because the spectral shaping of the input signal is adjusted according to the needs of the input signal to provide high masking at all times, and headroom is maintained throughout the frequency range because extreme preemphasis is used only when it is really needed.\nWhile the spectral compressor continuously adapts its characteristics to the specific nature of the program material, typically a significant overall preemphasis is needed during encoding to provide adequate masking. Instead of designing the spectral compressor to operate only as a preemphasis network, the BTSC system lets the spectral compressor, and hence the spectral expander, operate with a symmetric boost and cut, while obtaining the bias toward preemphasis from fixed preemphasis networks.\nThe spectral compander works in conjunction with the wideband compressor, thus the 2:1 wideband compressor senses the input signal level and adjusts a variable-gain stage to reduce the amplitude of large input signals and boost the amplitude of small input signals. This way, the amplitude dynamic range at the output of the wideband compressor is half that at the input.\nThe wideband compressor produces an output signal of relatively constant level. The choice of this level affects how well the noise is masked, higher levels being better, and how cleanly transients are reproduced, lower levels being better. The BTSC Standard has set this level at about 17 dB below 100% modulation at 300 Hz. This allows room for transients to overshoot in the compressor without causing excessive distortion through a complete companding cycle and ensures the program signal has a sufficiently high amplitude to mask the background noise.\nOn decode, the wideband expander restores the original program dynamics. During quiet passages, the expander attenuates the background noise, while during passages of significant signal amplitude level, the program itself masks the noise.\nAnother important concern of the BTSC design is to protect against large transients causing excessive modulation of the transmission system. In a peak-limited medium, such as the MTS transmission system, the peak excursion of the compressor output is an important parameter to control. This requirement can most easily be met by a clipper or limiter.\nClippers are relatively inaudible in operation provided that the clipper clips only transients that last under a few milliseconds. To accomplish this the unaffected-level point of the wideband compressor is set below 100% modulation point (Unaffected Level); and the wideband and spectral compressors, which precede the clipper, are set fast enough to allow only brief overloads to reach the clipper. Furthermore, the static preemphasis precedes the clipper, further reducing the transient overload duration. Aiming at a signal level for the output of the encoder that is substantially below 100% modulation, provides room for the typical peak excursions that occur when normal musical and other program transients are input to the compressor.\nAn important benefit of the Unaffected Level alignment is that the signal level broadcast over the stereo-subcarrier has its amplitude distribution between 10% to 30% modulation. This allows for large amplitudes in the monophonic carrier without exceeding the allowable modulation limits of monophonic plus stereophonic carriers. Since the stereo-subcarrier tends to average below 30% modulation, the mono carrier is allowed to stay around 70% modulation.\nTo protect the 15.734 KHz pilot signal and to prevent any spillover of signals into the sum channel, the difference channel must not contain information above about 15 KHz. This is accomplished by a lowpass filter in the encoder. For example, an 11-pole Cauer filter has been used for such application. This filtering process causes a phase shift in the difference channel (L-R signal) which must be compensated for in the sum (L+R)channel. This lowpass filter is located at the encoder output, but within the feedback loop, so that the RMS-level detectors in both the compressor and expander sense the same bandlimited signal. This is a limitation in prior art BTSC encoding, as it often requires high order filters and, as discussed below, often requires the use of matching filter components.\nThe ideal low-pass filter should pass the frequencies from 30 Hz to 15 KHz with minimal attenuation, have peak attenuation at 15.734 KHz, and maintain substantial attenuation above 15.734 KHz. Specific requirements for this filter can be found in the EIA document \"BTSC System Television Multichannel Sound Recommended Practices\". Note that the phase shift introduced by this filtering must be compensated for by an identical Phase shift in the sum (L+R) channel, or stereo separation will be significantly degraded. This is usually accomplished by using identical, matched filters in the L-R and L+R channels. This is a limitation in prior art BTSC encoding, since filter matching increases manufacturing costs, and errors in filter matching degrade performance.\nThe received demodulated difference (L-R) signal has high-frequency components caused by the other audio channels in the system, which could interfere with proper decoding if they are not attenuated before reaching the expander detectors. Therefore, filtering is necessary to prevent decoder mistracking. Pilot cancellation techniques will reduce the amount of 15.734 KHz in the L-R signal, and therefore reduce the required attenuation of the filter. There are many possible circuit variations which call for different alignments for this filter. For example, the filter may be located in the L-R signal path. However, this requires a compensating network, which in turn usually requires a matching filter to be located in the L+R signal path to maintain separation. Again, this exemplifies a limitation in prior art BTSC encoding, by requiring filter matching and additional hardware components.\nOST 60 describes a theoretically perfect compressor. Due to bandwidth limitations of presently available components, it may not be practical to conform perfectly to the ideal design of OST 60. This means that small deviations from perfect amplitude and phase response may exist in practical encoders, especially at the edges of the audio band. Such deviations may be compensated for by amplitude and phase errors introduced in the sum channel. Note that any deviations from perfect amplitude response in the encoder will be exaggerated by the expansion action of the decoder. For this reason, somewhat more amplitude response error must be placed in the sum channel than in the difference channel if the sum-channel response is to match that of the difference channel, including an encode/decode cycle. The amount of exaggeration of the error varies with frequency because the effective compression ratio varies with frequency. However, encoder phase errors should be compensated for with identical errors in the sum channel, since the decoder will not exaggerate phase errors.\nThe above discussion illustrates a significant limitation in prior art BTSC encoding. The hardware components introduce amplitude and phase errors, which errors, if not properly corrected or compensated for, will be magnified upon decoding; thus necessitating a frequency dependent overcompensation of the error. Furthermore, in prior art practice, it is an already accepted fact that due to bandwidth limitations of presently available components, it is not practical to conform perfectly to the ideal design of OST 60. Being unable to achieve the ideal requirements, the industry deviates from the ideal design during the signal encoding and compression. The signal decoding and decompression does not always compensate for such deviations, resulting in less than ideal noise reduction techniques.\nThe sum-channel compensation at the transmission end must correct only for phase and amplitude errors introduced by the encoder, while errors introduced by the decoder must be corrected only at the reception end. This allows freedom for future improvements in the state of the art. Furthermore, by compensating for transmission errors only in the transmitter, receiver manufacturers are free to build more nearly ideal receivers, rather than being forced to build errors into the receiver which compensate for errors in the transmitter. Also, by having the receiver end and the transmitter end compensate for their errors and limitations, this maintains receiver and transmitter independence and it does not impose a burden on either end to correct or compensate for errors or limitations of the other end.\nAs it is inherently typical with the adoption of a standard, the standard constrains the users to work and use that standard alone, thereby often stagnating the technological growth in the relevant field, or otherwise constraining the standard users and limiting the use of technology that may become superior by virtue of improvements in technology, or reductions in manufacturing costs. In the case of stereo-TV, because of compatibility concerns, a new and modem television set must be compatible with a 10 or 20 year old standard, effectively denying a consumer the enjoyment that could otherwise be available by improvements in technology.\nPrior art BTSC systems are effective in compressing and expanding the signal upon reception. However, prior art systems often are not very sensitive and responsive to transients where the signal amplitude quickly increases or decreases. Specifically, many expanders are unable to follow rapid envelope changes without producing undesirable acoustic effects, such as pumping and breathing. Many expanders try to overcome the above mentioned problems by slowing down its response to sudden amplitude changes. This, however, results in a response that is too slow to follow fast changing musical envelopes. Attempts to increase the response to said sudden amplitude changes often results in increased low frequency distortion and the above-mentioned pumping and breathing.\nOne existing system attempts to cure the above mentioned problems by varying the amplitude of the control signal as a function of the time derivative of the control signal. To achieve this, one system employs a lead-lag circuit including a diode. The diode is forward biased, resulting in fast reaction time to positive changing signals. However, the proportional derivative and lead-lag circuits are bi-polar and are unable to distinguish between positively increasing and negatively increasing signals. As a result, the response to negatively changing signals typically result in psychoacoustic distortion. In addition, the release behavior of the gain control module disclosed therein is non-linear which also affects the envelope and results in distortion. Furthermore, the circuit requires the selection of analog components to set a precise time-constant for a particular type of release, which time constant may be satisfactory for some types of programs, but not for others.\nAnother system attempts to overcome the diode limitation problems and non-linear circuit behavior. However, this is accomplished again via the use of diodes and other non-linear components, thereby requiring more complex circuit systems, extensive compensation and component selection.\nAccordingly, there is a need for a means to achieve BTSC encoding and decoding without having to contend with the non-linearity of components, independent of component tolerances and independent of variations in environmental conditions.\nAlso, there is a need to encode and decode a signal that is fully compatible and in compliance with current Multichannel Television Sound BTSC systems that does not deviate from the ideal low-pass filter recommended in the EIA document \"BTSC System Television Multichannel Sound Recommended Practices,\" that does not introduce undesirable phase shift or distortions, and that does not require the matching of filters or other components.\nAlso, there is a need for an encoder/decoder system, method and apparatus that exhibits greater signal to noise separation, improved performance characteristics throughout the spectrum of the signal, and improved response time and performance characteristics to sudden changes in the amplitude and/or frequency of the signal.\nThere is also a need to provide for the above-mentioned improvements at a reduced cost, while providing the service providers with means to adjust the encoder/decoder performance characteristics according to the type of program being transmitted, and preferably being adjustable on a real-time basis according to the instantaneous needs of the program.\nThere is also a need for an encoding/decoding system that can grow with and change with the changing needs in the market and advances in the technology, and is not bound or limited to perform according to a fixed standard.\nThere is further a need to provide for the above-mentioned improvements without sacrificing compatibility with current BTSC standards and without sacrificing performance on prior art television sets."}
{"text": "In recent years, there has been known a technique for detecting a lane mark such as a white line on a road where a vehicle travels by processing an image of the road acquired with an imaging means such as a CCD camera mounted on the vehicle and for performing steering control of the vehicle or providing a driver with information on the basis of information on a lane (traffic lane) along which the vehicle travels recognized from a detection result (refer to, for example, Japanese Patent Laid-Open No. Hei 11-147473/1999) (hereinafter, referred to as Patent Document 1).\nThe steering force assisting device in Patent Document 1 differentiates, with respect to a plurality of horizontal lines on the image of the road, luminance for each horizontal line from the left in the lateral direction, and extracts a point where luminance changes from dark to light (positive edge point) and a point where luminance changes from light to dark (negative edge point) on the basis of respective peak of the differential values. Then, a combination of edge points, in which the positive edge point and the negative edge point appear in alternate order on each horizontal line and in which the edge points are arranged at intervals that seem to be appropriate for a white line, is extracted as a white line candidate. Then, a white line is detected among the extracted white line candidates on the basis of the positions thereof on the image.\nOn the road, however, there may be a structure such as, for example, a guard rail besides the lane mark such as the white line. In this case, according to the technique in Patent Document 1, a portion of the guard rail is also extracted as edge points since the portion of the guard rail changes in light and dark on the acquired image. Therefore, the portion of the guard rail is also extracted as a white line candidate and may be detected as a white line. Particularly in the case where a structure such as a guard rail is installed on the road, a lane mark such as a white line defining the lane may be missing in some parts or temporarily invisible. In this case, the portion of the guard rail is more likely to be detected as a white line.\nMoreover, the structure such as a guard rail actually has a height and therefore a position where the guard rail is projected on the road surface is detected as a white line in the white line detection. Consequently, if recognizing the lane along which the vehicle travels from the position of the detected white line, the position and width of the lane would be different from that of the actual lane. Accordingly, the conventional device has a problem that the steering control of the vehicle based on the lane recognized in this manner causes inappropriate steering control of the vehicle."}
{"text": "Single crystal photovoltaic devices have been utilized for some time as sources of electrical power because they are inherently non-polluting, silent, and consume no expendable natural resources in their operation. However, the utility of such devices is limited by problems associated with the manufacture thereof. More particularly, single crystal materials (1) are difficult to produce in sizes substantially larger than several inches in diameter, (2) are thicker and heavier than their amorphous counterparts, and (3) are expensive and time consuming to fabricate.\nRecently, considerable efforts have been made to develop systems for depositing amorphous semiconductor materials, each of which can encompass relatively large areas, and which can be doped to form p-type and n-type materials for the production of p-i-n type photovoltaic devices which are, in operation, substantially equivalent to their crystalline counterparts. It is to be noted that the term \"amorphous,\" as used herein, includes all materials or alloys which have long range disorder, although they may have short or intermediate range order or even contain, at times, crystalline inclusions.\nIt is now possible to prepare amorphous silicon alloy material by glow discharge chemical vapor deposition. The silicon alloy material possesses (1) acceptably reduced concentrations of localized states in the energy gaps thereof, and (2) high quality electronic properties. Such techniques have been fully described in U.S. Pat. No. 4,226,898, entitled Amorphous Semiconductors Equivalent To Crystalline Semiconductors, issued to Stanford R. Ovshinsky and Arun Madan on Oct. 7, 1980, the disclosure of which is incorporated herein by reference. As disclosed in this patent, fluorine introduced into the discrete layers of amorphous silicon alloy material operates to substantially reduce the density of the localized states therein and facilitates the addition of other alloying materials, such as germanium.\nUnlike crystallize silicon which is limited to batch processing for the manufacture of solar cells, amorphous silicon alloys can be deposited in multiple layers over large area substrates to form solar cells in a high volume, continuous processing system. Such continuous processing systems are disclosed in U.S. Pat. No. 4,400,409 entitled \"A Method Of Making P-Doped Silicon Films And Devices Made Therefrom\"; U.S. Pat. No. 4,410,558 entitled \"Continuous Amorphous Solar Cell Production System\"; U.S. Pat. No. 4,438,723, entitled \"Multiple Chamber Deposition And Isolation System And Method\"; U.S. Pat. No. 4,492,181 entitled \"Method and Apparatus For Continuously Producing Tandem Amorphous Photovoltaic Cells\"; and U.S. Pat. No. 4,485,125 entitled \"Method and Apparatus For Continuously Producing Tandem Amorphous Photovoltaic Cells\", the disclosures of which are incorporated herein by reference. As disclosed in these patents, a substrate may be continuously advanced through a succession of interconnected deposition chambers, wherein each chamber is dedicated to the deposition of a specific semiconductor material. In making a photovoltaic device of p-i-n type configuration, the first chamber is dedicated for depositing a p-type layer of silicon alloy material, the second chamber is dedicated for depositing an intrinsic layer of amorphous silicon alloy material, and the third chamber is dedicated for depositing a layer of n-type silicon alloy material.\nSince each deposited layer of silicon alloy material, and especially the intrinsic layer, must be of high purity, (1) the deposition environment in the intrinsic deposition chamber is isolated, by specially designed gas gates, from the doping constituents introduced into the constituents into the intrinsic chamber, (2) the substrate is carefully cleansed prior to initiation of the deposition process to remove contaminants from the surface thereof, (3) all of the chambers which combine to form the deposition apparatus are sealed and leak checked to prevent the influx of environmental contaminants, (4) the deposition apparatus is evacuated and flushed with a sweep gas to remove contaminants from the interior walls thereof, and (5) only the purest reaction gases are employed from which to form the layers of silicon alloy material. In other words, every possible precaution is taken to insure that the sanctity of the vacuum envelope formed by the plurality of interconnected chambers of the deposition apparatus remains uncontaminated by impurities, regardless of origin.\nThe layers of silicon alloy thus deposited in the vacuum envelope of the deposition apparatus may be utilized to form a photovoltaic device including one or more p-i-n cells, one or more n-i-p cells, Schottky barrier devices, photodiodes, phototransistors, or the like. Additionally, by making multiple passes through the succession of deposition chambers, or by providing an additional array of deposition chambers, multiple stacked cells of various configurations may be obtained.\nMost photovoltaic devices, having either single cell or multiple cell structures, preferably include a light reflecting back reflector for increasing the percentage of incident light reflected from the substrate back through the active layer of silicon alloy material of the cells. It should be obvious that the use of a back reflector increases the amount of light which passes through the active layer of silicon alloy material, thus increasing the amount of incident light which is converted to electricity, and increasing the operational efficiency of the photovoltaic device. However, all layers, other than the photoactive layers of silicon alloy material, deposited atop the light incident surface of the substrate must be substantially transparent (to 35 nanometers to one micron wavelength light) so as to pass a high percentage of incident light from the anti-reflective coating atop the photovoltaic cell to the highly reflective surface of the back reflector from which it is redirected through and absorbed by the photoactive layer of silicon alloy material.\nThe back reflector may be formed atop the deposition surface of the substrate if an opaque substrate is employed. The back reflector may be either specular or diffuse. With either type of back reflector, light which has initially passed through the layers of silicon alloy material from which the photovoltaic device is fabricated without being absorbed, is redirected by the highly reflective material of the back reflector to pass, once again, through the photoactive layers. The additional pass results in increased photon absorption and charge carrier generation, thereby providing increased short circuit current.\nThe aforedescribed thin film amorphous silicon alloy materials offer several distinct advantages over crystalline materials, insofar as they can be easily and economically fabricated to cover large areas by the newly developed mass production processes. However, in the fabrication of said silicon alloy materials by the aforementioned processes, the presence of current-shunting defects has been noted. These defects have (1) seriously impaired the performance of the photovoltaic devices fabricated therefrom and (2) detrimentally affected production yield. These process-related defects are thought to either (1) be present in the morphology of the substrate electrode, or (2) develop during the deposition of the multiple layers of silicon alloy material. The instant invention is directed toward eliminating, or at least substantially reducing, the effects of these current-shunting defects.\nThe most important of these defects may be characterized as \"shunts,\" \"short-circuits,\" defect regions, or low resistance current paths. Before the suspected causes of these defects are explained, it is helpful to note the thicknesses of the deposited of silicon alloy material layers. In a typical p-i-n type photovoltaic device, the \"p\" layer is on the order of 75-200 Angstroms thick, the intrinsic layer is on the order of 3,500 Angstroms thick, and the \"n\" layer is on the order of 75-200 Angstroms thick, thereby providing a total semiconductor body thickness of only about 4,000 Angstroms. It should therefore be appreciated that irregularities disposed on the deposition surface, however small, may not be fully covered by the conformally deposited layers of silicon alloy material.\nShunt defects are present when one or more low resistance current paths develop between the upper and lower electrodes of the photovoltaic device. Under operating conditions, a photovoltaic device in which a shunt defect has developed, exhibits either (1) low power output, since electrical current collected at the electrodes flows through the defect region (the path of least resistance) in preference to an external load, or (2) complete failure where sufficient current is shunted through the defect region to \"burn out\" the device.\nWhile shunt-type defects always deleteriously affect the performance of photovoltaic devices, their effect is greatest when the devices in which they are incorporated are operated under relatively low illumination such as room light, vis-a-vis, high intensity illumination such as AM-1. Under room light illumination, the load resistance of the cell (i.e., the resistance under which the cell is designed to operate most efficiently) is comparable to the shunt resistance (i.e., the internal resistance imposed by the defect region), whereas under AM-1 illumination, the load resistance is much lower by comparison. This occurs because, in a photovoltaic device, photogenerated current increases linearly with increasing illumination, while the resulting voltage increases exponentially. In other words, voltage attains a relatively high value under low illumination, the value increasing only slightly as the intensity of the illumination is increased. Therefore, under low intensity illumination the relatively high voltage potential present preferentially drives the relatively small number of photogenerated current carriers through the path of least resistance, i.e., the low resistance defect regions. In contrast thereto, under high intensity illumination, a large number of current carriers are present and are driven by a potential of about the same magnitude as the potential which exists under low illumination. This larger number of current carriers compete for a limited number of least resistance paths (through the defect regions). The result is that under high intensity illumination, while more power may be lost to the defect region, the power lost is a smaller percentage of the total photogenerated power than under low intensity illumination.\nDefects or defect regions, the terms being interchangeably used herein, are not limited to \"overt\" or \"patent\" short circuit current paths. In some cases, the adverse effects of a defect are latent and do not immediately manifest themselves. Latent defects can give rise to what will be referred to hereinafter as an \"operational mode failure,\" wherein a photovoltaic device, initially exhibiting satisfactory electrical performance, suddenly fails. The failures will be referred to herein as operational mode failures regardless of whether the device was previously connected to a load for the photogeneration of power, it only being necessary that the device was, at some time subjected to illumination, thereby initiating the photogeneration of charge carriers. It is believed the shunt defects, both latent and patent, arise from one or more irregularities in the (1) morphology of the exposed surface of the substrate material, or (2) in the growth of the layers of silicon alloy material thereupon.\nThe first, and perhaps most important, source of the defects, i.e., the aforementioned morphological irregularities existing on the deposition surface of the substrate material, will now be discussed. Even though the highest quality, for example, stainless steel is employed as the substrate upon which the layers of silicon alloy material are successively deposited, it has been calculated that from 10,000 to 100,000 irregularities per square centimeter are present on the deposition surface thereof. Such irregularities take the form of projections, craters, or other deviations from a smooth finish and may be under a micron in (1) depth below the surface, (2) height above the surface, or (3) diameter. Regardless of their configuration or size, the defects may establish a low resistance current path through the layers of silicon alloy material, thereby effectively short-circuiting the two electrodes. This may occur in numerous ways. For instance, a spike projecting from the surface of the substrate electrode may be of too great a height to be covered by the subsequent deposition of the thin film layers of silicon alloy material, and therefore, be in direct electrical contact with the other light incident electrode when that upper electrode is deposited atop the layers of silicon alloy material. Likewise, a crater formed in the surface of the substrate electrode may be of too large a size to be filled by the subsequent deposition of the silicon alloy material. In such an instance: (1) electrical current may bridge the gap which exists between the electrodes, or (2) during actual use (the photoinduced generation of electrical current) of the photovoltaic device, the material of one of the electrodes may, under the influence of the electrical field, migrate toward and contact the other of the electrodes, and thereby pass electrical current therebetween. It is also possible that in some cases the layers of silicon alloy material deposited onto the substrate include regions of irregular composition which can provide low resistance paths for the flow of electrical current between the electrodes of the photovoltaic device.\nFurther, despite all the previously described efforts to maintain the vacuum envelope free of external contaminants, dust or other particulate matter, which somehow either (1) invades the vacuum envelope during the deposition of the silicon alloy material, or (2) forms as a by-product of the deposition process, may be deposited over the substrate electrode along with the silicon alloy material. The contaminants interfere with the uniform deposition of the layers of silicon alloy material and may establish low resistance current paths therethrough.\nAdditionally, it is suspected that in some case, the silicon alloy material may form micro-craters or micro-projections during the deposition thereof, even absent the presence of contaminants or pollutants from external sources. Such morphological deviations from a perfectly smooth and even surface means that the substrate is covered by silicon alloy material either (1) in an \"ultra thin layer\" (consider again that the total thickness of all layers is only on the order of 4,000 Angstroms and any reduction in coverage is indeed an ultra thin layer) or (2) not at all. Obviously, when the upper electrode material (typically a conductive transparent oxide) is deposited across the entire surface of the body of silicon alloy material, the defect regions cause the low resistance current path to develop, and electrical current is shunted therethrough. In still other cases involving defect regions, the presence of such defect regions is only detectable due to their deleterious effect upon the electrical and photoelectric properties of the resultant photovoltaic device. Finally, note that the defects described hereinabove may not be sufficiently severe to divert all electrical current through the low resistance path. However, the diversion or shunting of any current therethrough represents a loss in operational efficiency of the photovoltaic device and should therefore be eliminated. Moreover, the shunting of even small amounts of current through each of thousands of defect regions may combine to cause major losses in efficiency. Based upon the foregoing, it should be apparent to one ordinarily skilled in the art that a reduction in current flow through these defects and defect regions is critical to the fabrication of high-yield, high efficiency, large area thin film photovoltaic devices.\nSeveral approaches with which to cope with gross amplifications of this problem have been implemented by the instant inventors and their colleagues. As described in U.S. Pat. No. 4,451,970 of Masatsugu Izu and Vincent Cannella, entitled \"System and Method for Eliminating Short Circuit Current Paths In Photovoltaic Devices\", the disclosure of which is incorporated herein by reference, the shunting of current through defect regions is treated by substantially eliminating those defect regions as an operative area of the semiconductor device. This is accomplished in an electrolytic process where electrode material is removed from the periphery of the defect site, effectively isolating the defect regions and preventing the flow of electrical current therethrough. However, the process described in the '970 patent is current dependent, i.e., the greater the amount of current flowing through a particular area of the device, such as a defect region, the greater the amount of electrode material (in the preferred embodiment indium tin oxide) removed. Consequently, said short circuit eliminating process performs admirably in removing the electrode material from the periphery of a large defect, and thereby preventing all current flow therethrough. However, since the method is current dependent, it is not as successful in eliminating the flow of current between the electrodes in the thousands of defect regions which are relatively small. And as previously mentioned, since a great many relatively small current shunting paths, taken in toto, divert a substantial amount of current from its desired path of travel, the low resistance current paths created by such small defect regions must also be eliminated or at least substantially reduced. Further, the electrolytic process described in the '970 patent neither detects nor helps in preventing the formation of current-shunting paths in the case of operational mode failures.\nIn U.S. Pat. No. 4,419,530 of Prem Nath, entitled \"Improved Solar Cell and Method For Producing Same\", the disclosure of which is incorporated herein by reference, there is described a method for electrically isolating small area segments of amorphous, thin film, large area photovoltaic devices. This isolation of defects is accomplished by (1) dividing the large area device into a plurality of small area segments, (2) testing the small area segments for electrical operability, and (3) electrically connecting only those small area segments of the large area device exhibiting a predetermined level of electrical operability, whereby a large area photovoltaic device comprising only electrically operative small area segments is formed. While this method effectively reduces or eliminates the effect of defects, it is not completely satisfactory for several reasons.\nThe step of dividing the large area solar cell into electrically isolated small area segments requires several production steps and also reduces the total area of the solar cell that is available for producing electrical energy. Further, the method can be time and cost intensive since the electrical output of each isolated portion must be tested and separate electrical connections must be made to provide electrical contact to each small area segment. Also, since entire small area segments are effectively eliminated from the final cell if they manifest a defect, proportional losses of efficiency are greater than they would be if only the precise area of the particular defect were eliminated. In addition, it is possible that defects (shorts) in a solar cell can develop after the cell has been in use, and the concept of dividing the surface of the large area cell into small area segments does not correct this type of latent defect.\nFurther, both of the foregoing patent applications relate to \"after market\" techniques which are applicable to (1) isolate only gross defect-containing regions, and (2) prevent any and all current flow through those defect containing regions. Accordingly, a need still exists for a photovoltaic device which substantially eliminates the deleterious effects of shunts and other defects, both large and small, whatever their origin, without operatively removing large portions of the active body of silicon alloy material or the electrode disposed thereupon, while maintaining an acceptable level of current flow across the entire surface of the device.\nAnother such method and device is disclosed in\nU.S. Pat. No. 4,590,327, of Nath, et al, entitled \"Photovoltaic Device and Method\", the disclosure which is incorporated herein by reference. Disclosed therein are several configurations of current collecting bus grid structures for photovoltaic devices designed to minimize the effects of shorts, shunts, and other defects upon the performance of the devices. The structures for reducing the degrading effect of a low resistance current disclosed therein can be (1) a specifically designed layer of transparent electrically conductive material, (2) electrical isolation of the current carrying portions of the bus grid structure from the conductive layer, (3) resistive connections of the current collecting fingers to the remainder of the bus grid structure, (4) a buffered bus grid structure or a fuse-type connection, or (5) the grid or current collecting fingers which terminate when the current reaches a predetermined amount. While these solutions effectively reduce or eliminate the effect of defects, they are not completely satisfactory for several reasons. The specifically designed layers of transparent electrically conductive material are difficult to manufacture. Further, devices having terminable or otherwise destructible current collecting fingers reduce surface area otherwise available to photogenerate current as the fingers are terminated.\nA differently configured photovoltaic device and method for eliminating the problems of shorts and shunts is disclosed in U.S. Pat. No. 4,633,034 of Nath, et al, entitled \"Photovoltaic Device and Method\", the disclosure of which is incorporated herein by reference. That application concerns a photovoltaic device having the current collecting bus grid structure disposed above the upper, transparent conductive electrode thereof. A pattern of electrical current flow restricting material is disposed between the layer of transparent conductive material and the substrate of the device, corresponding to, and disposed beneath, at least a portion of the superjacent current collection means. The pattern restricts the flow of electrical current between the portions of the current collection means and the substrate positioned beneath the body of amorphous silicon alloy material. However, it has been determined from actual manufacturing experience that such an internal current restricting material must be aligned with the current collection means to operate effectively, which requires exacting manufacturing techniques.\nAccordingly, there exists a need for a defect tolerant photovoltaic device manufactured by a process which does not require exacting manufacturing techniques such as those required to align subjacent components thereof or to prepare specifically designed layers of transparent electrically conductive material.\nAdditionally, it must be pointed out that as much as 10% of the surface of a photovoltaic cell is occupied by a bus grid structure. Any reduction in this area, dedicated for use by the bus grid structure, would be welcomed because radiation incident thereupon is lost. More particularly, due to \"shadowing,\" radiation incident upon said structure cannot pass to the photoactive region of the cell for the photogeneration of electron-hole pairs therein. If the thin grid fingers could be made longer, at least some of the wide bus bars could be eliminated, thereby increasing the active area of the photovoltaic cell available to photogenerate electron-hole pairs. However, in order to be able to elongate the grid fingers, the cell must be made more tolerant of short circuit defects than is currently the case.\nDisclosed herein as the innovative features of the instant invention are configurations of photovoltaic devices which exhibit a high degree of operational tolerance to defects therein, and which utilize a novel current collecting bus grid structure disposed above the transparent conductive top electrode of the photovoltaic device, which bus grid structure can be optimized to maximize the active area of a photovoltaic device."}
{"text": "The present invention relates broadly to home appliances for cooking and, more particularly, to a gas range having a supplemental primary air supply.\nGas ranges provide an excellent baking and broiling environment for preparing various dishes in a highly effective manner. Gas ovens burn hot and steady and provide excellent temperature control as well as an instant on-and-off operation with little residual heat. For oven heating, gas ovens include bake elements and broil elements similar to their electric counterparts, but with different structures based on the need to supply the gas heating elements with a combustible fuel and air mixture and the ability to produce discrete flames for even oven-cavity heating.\nGas broiler elements are used when the oven is to be taken to an elevated temperature for broil-type cooking and due to the location of the broil element near the top of the oven cavity, the space between the oven cavity and an associated cooktop can reach undesirable temperatures. Therefore, in order to provide better control of the cooktop temperature as well as protect the electronics of nearby control systems, gas ovens can include a ventilation channel that can provide a cooling airstream intermediate the oven cavity and the cooktop. Usually such a ventilation channel includes a fan for forced-air ventilation of the channel intermediate the cooktop and the oven cavity.\nThe gas burner is fueled by a mixture of gas injected by a gas nozzle and ambient air from behind the oven cavity. The air-to-gas ratio is controlled generally to provide the most efficient burning situation while providing the necessary energy input to elevate the oven to cooking temperatures, or in the case of pyrolytic cleaning, to cleaning temperatures.\nDuring times of elevated temperatures within the oven cavity and when the fan is in operation, the fan intake tends to draw primary air away from the broil burner intake and the burning efficiency of the broil burner is reduced due to air deficiency. It therefore becomes desirable to enhance the efficiency of the gas broil burner by providing more primary air to the broil burner intake during such periods of air deficiency."}
{"text": "Diesel powered systems such as, but not limited to, off-road vehicles, marine diesel powered propulsion plants, stationary diesel powered system and rail vehicle systems, or trains, are usually powered by a diesel power unit. With respect to rail vehicle systems, the diesel power unit is part of at least one locomotive and the train further includes a plurality of rail cars, such as freight cars. Locomotives are complex systems with numerous subsystems, with each subsystem being interdependent on other subsystems.\nDiesel fueled power units typically require cooling systems to limit the temperatures of various engine components. Internal combustion engines have internal cooling passages for the circulation of coolant to remove heat energy from the engine components. Lubricating oil which is circulated throughout the engine to reduce friction will also absorb heat and, therefore, will also require cooling to avoid reaching temperatures that would detrimentally affect its lubricity. Diesel engines often utilize turbochargers to increase power by compressing the intake combustion air to a higher density. Such compression results in the heating of the combustion air, which must then be cooled prior to its use to enable the engine to have high volumetric efficiency and low emissions of exhaust pollutants. For mobile applications such as rail locomotives, the only readily available heat sink is the surrounding ambient air. It is known to utilize a pumped cooling medium, such as water, to transport heat to finned radiator tubes. The radiator tubes then transfer the heat to the ambient air, often using forced convection provided by fans.\nIt is often desirable to maintain an internal combustion engine and its associated intake combustion air at multiple different temperatures in order to optimize the performance of the engine. Consequently, coolant at one temperature may be provided to the cylinder jackets of a turbocharged diesel engine and coolant at a lower temperature may be provided to an intercooler for cooling the compressed combustion air. Such a system may use a single pump, heat exchanger, and temperature control valve to accomplish the dual cooling objectives.\nOther turbocharged diesel engine cooling schemes may use a subcooler in addition to a radiator. The subcooler is typically located upstream of the radiator in a flow of cooling ambient air. For locomotive applications, ambient air flowing through the radiators is normally provided by a multi-speed fan, since the radiators are positioned on the roof of the locomotive. The use of a subcooler provides a greater temperature difference capability between the temperature of the engine and the temperature of the combustion air.\nDiesel engines may also use cooled fluid and/or compressed fluid cooling systems to cool lubricant fluids, such as engine lubrication oil. For example, heated engine oil coming from the engine may be passed through an oil cooler having an oil/coolant heat exchanger supplied with a coolant at a lower temperature than the oil to transfer heat from the oil to the coolant. Typically, the coolant includes water and/or a water antifreeze mixture cooled in the radiator and/or one or more associated sub-coolers. Additional oil cooling may be provided using a compressed coolant cooling scheme.\nConventional turbocharged diesel powered locomotives may use a cooling circuit that selectively uses a sub-cooler radiator section to provide sub cooled water or an antifreeze/water mixture to an oil cooler and/or the turbocharger compressor intercooler. At lower ambient temperatures and or low load conditions, the conventional locomotive cooling circuit may provide coolant to both the oil cooler and the intercooler to maintain a desired engine oil temperature and a desired emission level produced by the engine. During operation in high ambient temperatures and/or at high engine loading, it may not be possible to provide sufficient cooling capacity from the coolant in the circuit to cool the engine oil without sacrificing cooling of another cooled component, such as the intercooler. Consequently, at higher ambient temperatures and/or higher loads, coolant supplied to the intercooler may be redirected to the oil cooler to preferentially cool the engine oil at the expense of higher pollutant emission. For example, U.S. Pat. No. 7,131,403 describes one such cooling system for maintaining sufficient coolant flows to the intercooler and oil cooler for providing lower NOx emissions at ambient temperatures up to about 80° Farenheit (F.). Above 80° F. ambient, the cooling system may be configured for routing coolant flows so as to protect the engine by not exceeding oil temperature limits, but at the expense of producing higher NOx emissions.\nOwners and/or operators of locomotives, off-road vehicles, marine diesel powered propulsion plants, and/or stationary diesel powered systems desire to optimize cooling of diesel fueled power generation units used in these applications over a wide range of operating conditions and ambient environments."}
{"text": "1. Field of the Invention\nThe present invention relates to a lamp assembly, and particularly to a light emitting diode lamp assembly employed in a direct type backlight module of a liquid crystal display.\n2. Discussion of the Related Art\nTypically, light emitting diodes (LEDs) are preferred over other types of light sources because LEDs exhibits low energy consumption, long service life, and other advantages. Therefore, LEDs are widely used as light sources.\nFIG. 1 illustrates a typical LED lamp assembly 10 using LEDs as a light source. The LED lamp assembly 10 includes a housing 11, a printed circuit board 12, a plurality of side-lighting type LEDs 13, a light reflective module 14, a plurality of circular reflective layers 15 and a lamp cover 16. The housing 11 is an elongated, hollow structure having an opening 112. The LEDs 13 are arranged apart and electrically connected to the printed circuit board 12. The printed circuit board 12 with the LEDs 13 is disposed on a bottom surface of the housing 11. The light reflective module 14 includes a rectangular bottom reflective plate 144 and four connecting sidewalls 142 extending from a periphery of the bottom reflective plate 144. The bottom reflective plate 144 defines a plurality of through holes 146, configured for allowing the light-emitting portions of the LEDs 13 to pass through. The light reflective module 14 can be correspondingly mounted into the housing 11 via the opening 112. The circular reflective layers 15 are positioned at the tops of the LEDs 13 respectively. The lamp cover 16 is fixed on the opening 112 of the housing 11. Light from the light-emitted portions of the LEDs 13 is substantially reflected at the sidewalls 142 and the bottom reflective plate 144, finally outputted from the lamp cover 16. With the help of the light reflective module 14, an efficiency of utilization of light energy of the LED lamp assembly 10 is increased.\nNevertheless, the brightness above the LEDs 13 of the LED lamp assembly 10 is decreased due to the light reflection off the circular reflective layers 15, and a plurality of dark areas between the two adjacent LEDs 13 still occur. Accordingly a uniform brightness of the LED lamp assembly 10 is low.\nWhat is needed, therefore, is an LED lamp assembly that overcomes the above mentioned disadvantages."}
{"text": "Sound sources are generally devices that generate acoustic energy. One use of sound sources is in marine seismic surveying in which the sound sources may be employed to generate acoustic energy that travels downwardly through water and into subsurface formations. After interacting with the subsurface formations, some of the acoustic energy may be returned toward the water surface and detected by specialized sensors, such as hydrophones or geophones. The detected energy may be used to infer certain properties of the subsurface formations, such as structure, mineral composition and fluid content, thereby providing information useful in the recovery of hydrocarbons.\nOne type of sound source that may be used in marine seismic surveying may be a marine acoustic vibrator. Increasing the acoustic output of marine acoustic vibrators may require an increase in input power that may result in an increase in heat generation. “Moving-coil,” also referred to as “voice-coil” transducers, may be employed in marine acoustic vibrators, for example, to provide the ability to generate very large acoustic amplitude. However, while some marine acoustic vibrators may exhibit comparatively high sound pressure levels versus previous sound sources, these marine acoustic vibrators may exhibit various challenges with thermal management, for example. One such challenge may include design of the coil element, which should be able to withstand high power densities while maintaining contact with the acoustic load (e.g., vibrating shell, etc.). These constraints may result in induced thermal strain within the marine acoustic vibrator and/or within the coil element itself."}
{"text": "This invention relates to the field of resonators and filters and more specifically to switchable and tunable coplanar waveguide resonators and filters using micro-electro-mechanical switches.\nA need exists for switchable and tunable filters for both wideband and multiband communication systems that are small in size, inexpensive and easy to manufacture. Prior art switchable/tunable filters use resonant ring arrangements with diodes used as switches to select filter response. These diode switches tend to be large in size and expensive. To help alleviate this problem, attempts have been made to manufacture filters using a micro-electro-mechanical system method. This leads to a switchable filter system but requires two different filter structures to be built with a diode for use as a switch to switch the signal path from one filter structure to another. Drawbacks to this design include that the process used to make these filters is complicated, that the resultant filter structure is large, and that there is interference between the two filter structures.\nAnother approach is to use coplanar waveguide filters. Coplanar waveguide filters are manufactured by using a substrate covered with a metal layer. The metal is etched to layout various filter configurations. While small filters can be manufactured in this manner, a switchable filter system still requires two different filters connected by a diode switch. This results in a large structure.\nWhat is needed is a filter that can combine both coplanar waveguide filters with micro-electro-mechanical switches."}
{"text": "Communication networks are in wide use in many technological fields including distributed computing, data exchange and telecommunication applications. Communication networks generally include a plurality of nodes, such as bridges, LAN switches, routers, cross-connections and telephone switches. The networks further include communication links, such as cables, point-to-point radio connections and optical fibers, which connect the nodes. The networks also include ports, generally within some of the nodes, for attaching external devices such as computers, terminals, handsets, and multiplexers referred to as end-points (or hosts).\nNetworks are becoming increasingly complex, especially due to their increasing speeds of operation, the number of units interconnected by a network and the formation of large networks from different types of sub-networks. In addition, the networks may transmit concurrently various types of data such as text, voice, video and other multimedia files. In order to allow for these different types of data, some networks are designed to provide different amounts of bandwidth and different levels of quality of service.\nA major issue with newly deployed and existing communication networks is testing and trouble-shooting, i.e., checking whether the network is operating according to its specifications and, if not, determining the cause of the network's inadequate performance (for example, the identity of a faulty unit).\nSimulators such as \"BONeS,\" produced by Cadence, San Jose, Calif., and \"OPNET,\" produced by MIL3, Washington, D.C., allow creation of models of a network, and subsequent testing of the network based on these models. An operator provides the simulator with a map of the network, which includes its components, specifications and expected data traffic patterns. The simulator provides performance estimations of the entire network, together with performance estimations of the network under various constraints, such as a non-operating node or link. Such performance estimation under constraints is referred to as \"what-if\" analysis. However, simulators do not usually cover all aspects of the simulated networks and are limited in the network size which may be simulated. In addition, simulators are able to detect only design faults, not operational faults.\nApplication performance measurement tools, such as \"Chariot,\" produced by Ganymede, Research Triangle Park, N.C., and \"Webload\" and \"Webexam,\" produced by Radview, Tel Aviv, Israel, evaluate the performance of existing or new applications as they are introduced into the network. However, they do not test the network itself independent of specific applications. In addition, they do not provide \"what-if\" analysis capabilities.\nDedicated point-to-point testing equipment is a commonly-used network testing tool. Such equipment is described, for example in U.S. Pat. No. 5,477,531, which is incorporated herein by reference. Usually, dedicated point-to-point testing equipment requires two users who coordinate their operations in order to identify a misbehaving component of the network. Moreover, to test a large network, the testing equipment must be moved between many ports of the network.\nPassive network monitoring devices based on network management software, such as \"HP-OpenView,\" produced by Hewlett Packard, Palo Alto, Calif., and \"Soltice Enterprise Manager,\" produced by Sun Microsystems, Mountain View, Calif., are also in very common use. These devices are used in network troubleshooting, but they are passive and cannot test a network which is not in use."}
{"text": "Mobile communications have changed the way people communicate and mobile phones have been transformed from a luxury item to an essential part of every day life. The use of mobile phones is today dictated by social situations, rather than hampered by location or technology. While voice connections fulfill the basic need to communicate, and mobile voice connections continue to filter even further into the fabric of every day life, the mobile Internet is the next step in the mobile communication revolution. The mobile Internet is poised to become a common source of everyday information, and easy, versatile mobile access to this data will be taken for granted.\nThird generation (3G) cellular networks have been specifically designed to fulfill these future demands of the mobile Internet. As these services grow in popularity and usage, factors such as cost efficient optimization of network capacity and quality of service (QoS) will become even more essential to cellular operators than it is today. These factors may be achieved with careful network planning and operation, improvements in transmission methods, and advances in receiver techniques. To this end, carriers need technologies that will allow them to increase downlink throughput and, in turn, offer advanced QoS capabilities and speeds that rival those delivered by cable modem and/or DSL service providers.\nIn order to meet these demands, communication systems using multiple antennas at both the transmitter and the receiver have recently received increased attention due to their promise of providing significant capacity increase in a wireless fading environment. These multi-antenna configurations, also known as smart antenna techniques, may be utilized to mitigate the negative effects of multipath and/or signal interference on signal reception. It is anticipated that smart antenna techniques may be increasingly utilized both in connection with the deployment of base station infrastructure and mobile subscriber units in cellular systems to address the increasing capacity demands being placed on those systems. These demands arise, in part, from a shift underway from current voice-based services to next-generation wireless multimedia services that provide voice, video, and data communication.\nThe utilization of multiple transmit and/or receive antennas is designed to introduce a diversity gain and to increase the degrees of freedom to suppress interference generated within the signal reception process. Diversity gains improve system performance by increasing received signal-to-noise ratio and by stabilizing the transmission link. On the other hand, more degrees of freedom may allow multiple transmissions that may provide more robustness against signal interference, and/or by permitting greater frequency reuse for higher capacity. In communication systems that incorporate multi-antenna receivers, a set of M receive antennas may be utilized to null the effect of (M−1) interferers, for example. Accordingly, N signals may be simultaneously transmitted in the same bandwidth using N transmit antennas, with the transmitted signal then being separated into N respective signals by way of a set of N antennas deployed at the receiver. Systems that utilize multiple transmit and receive antennas may be referred to as multiple-input multiple-output (MIMO) systems. One attractive aspect of multi-antenna systems, in particular MIMO systems, is the significant increase in system capacity that may be achieved by utilizing these transmission configurations. For a fixed overall transmitted power, the capacity offered by a MIMO configuration may scale with the increased signal-to-noise ratio (SNR). For example, in the case of fading multipath channels, a MIMO configuration may increase system capacity by nearly M additional bits/cycle for each 3-dB increase in SNR.\nThe widespread deployment of multi-antenna systems in wireless communications has been limited by the increased cost that results from increased size, complexity, and power consumption. This poses problems for wireless system designs and applications. As a result, some initial work on multiple antenna systems may be focused on systems that support single user point-to-point links. However, the use of multi-antenna techniques for a multiuser environment to improve total throughput remains a challenge.\nFurther limitations and disadvantages of conventional and traditional approaches will become apparent to one of skill in the art, through comparison of such systems with some aspects of the present invention as set forth in the remainder of the present application with reference to the drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a cable-type steering device in which a steering wheel and a steering gear box are connected to each other by flexible cables such as Bowden cables.\n2. Related Art\nA conventional steering device for a vehicle is designed such that a steering shaft having a steering wheel at its upper end is connected at its lower end to a steering gear box, so that a steering torque inputted to the steering wheel is transmitted through the steering shaft to a rack and pinion mechanism provided within the steering gear box.\nHowever, if the steering wheel and the steering gear box are connected to each other using the steering shaft, the following problem is encountered: it is difficult to freely select the position of the steering wheel relative to the position of the steering gear box. For this reason, not only the degree of freedom of the design is limited, but the steering gear box also cannot be commonly used in both of a right-side driven vehicle and a left-side driven vehicle. Moreover, there is another problem that the vibration received from a road surface through a tire and the vibration of an engine are transmitted through the steering shaft to the steering wheel, so that such vibrations detract from the calmness within a vehicle compartment and thus the riding comfort.\nTherefore, in place of the conventional steering shaft, a cable-type steering device has been proposed which uses a flexible transmitting means such as Bowden cables (for example, see Japanese Patent Application Laid-open No.8-2431).\nThe use of a cable-type steering device allows the position of the steering wheel relative to the position of the steering gear box to be freely selected. Additionally, the vibration of the steering gear box has difficulty being transmitted to the steering wheel. Thus, the above-described problems can be eliminated.\nWhen the cable-type steering device is combined with a power steering device, the driver's steering operation can be assisted by detecting a steering torque inputted to the steering wheel to drive an actuator for the power steering device. In this case, if a steering torque detecting means is provided in the steering gear box, the following problem is encountered: it is impossible to detect a steering torque including a friction of cables interposed between the steering wheel and the steering torque detecting means. For this reason, an assisting force offsetting the friction cannot be generated in the actuator, resulting in a reduced steering feeling."}
{"text": "This invention relates to a process for making silicon carbide whiskers and is particularly concerned with a process for providing relatively high yields of beta silicon carbide whiskers while controlling the size and shape of the whiskers.\nSilicon carbide is a high strength ceramic material which has good chemical stability and excellent oxidation resistance to high temperatures due mainly to covalent bonding and a crystal structure related to that of diamond. Silicon carbide whiskers are needle-shaped single crystals of silicon carbide having an aspect ratio, i.e., a length-to-diameter ratio, greater than about 3 and a typical diameter between about 0.1 and 10 microns. The high aspect ratio of whiskers makes them a much more effective reinforcement in composites, especially ceramic matrix composites, than silicon carbide particulates. Also, whiskers exhibit much higher mechanical strength than silicon carbide fibers, which are either polycrystalline or amorphous forms of silicon carbide that typically have a diameter greater than 10 microns. These superior properties of silicon carbide whiskers have led to their use as a reinforcing material for ceramics, metals, polymers and glass composites. Silicon carbide whiskers are particularly suited for use in the reinforcement of all types of engineering ceramics including gas turbine ceramics, automotive ceramics and ceramic cutting tools. It is estimated that the market for the use of whiskers as reinforcements for engineering ceramics alone will be several hundred million dollars per year by the year 2000.\nWhen silicon carbide whiskers are used as reinforcement in composites, it is preferable that the size and shape of the whiskers meet specified requirements. For example, in some instances it may be desirable that the diameter of the whiskers be less than one micron while in other cases whisker diameters of several microns are preferred. Also, straight whiskers are normally used to reinforce ceramic matrix composites because such reinforced composites have a higher fracture toughness than those reinforced with curly whiskers. On the other hand, curly whiskers may be more effective than straight whiskers for other reinforcing applications. For these and other reasons, it is desirable that the size and shape of the whiskers be controlled during synthesis.\nThe current dominant technology for producing silicon carbide whiskers is the carbothermal reduction of silica by carbon in rice hulls. Rice hulls are composed of about 15 to 20 percent ash that is primarily silica. Thus, when the rice hulls are retorted or pyrolyzed in an inert atmosphere at a temperature in the vicinity of 1800.degree.C. while removing gases and other vapors as they form, carbon in the rice hulls reacts with silica to form silicon carbide in accordance with the overall reaction of EQU 3C+SiO.sub.2 =SiC+2CO(g).\nThe major problem with using rice hulls to synthesize silicon carbide whiskers is that the resultant product normally contains only between about 10 and 20 weight percent silicon carbide whiskers with the remainder being silicon carbide in the form of particulates, unreacted silica and unreacted carbon. Higher yields of silicon carbide whiskers are usually not possible because the chemical composition of the rice hulls and the degree of mixing of the carbon and silica therein are set by nature and cannot readily be varied. Thus, it is difficult to obtain an intimate mixture of carbon and silica that is sufficiently porous to allow carbon monoxide gas to escape and thereby drive the overall reaction of carbon with silica to form silicon carbide to completion while allowing space for silicon carbide whiskers to grow.\nIt can be seen from the above that future processes for producing silicon carbide whiskers should give a relatively high yield of whiskers having desired sizes and shapes which can be tailored according to the particular use desired for the whiskers."}
{"text": "Automotive trim parts, particularly body side moldings of various sorts have been known for many years. Such trim parts have both decorative and protective functions. Typically, there are at least three pieces horizontally aligned and attached on the side of a vehicle: a piece for the front fender; a piece for each door; and a piece for the rear fender.\nOne known type of molding is made of a polymer coextruded onto a metal carrier strip. The metal is typically aluminum, stainless steel or bimetal and a portion of the metal may be left exposed to enhance the appearance of the molding. The underlying metal carrier generally has an outwardly convex cross-section and its longitudinal edges are bent inwardly to form lips running the length of the piece. This arrangement thus yields a hollow molding which is generally fastened to a vehicle by clipping of the lips to the vehicle body. This type of hollow molding requires a member to be fastened at each of its ends. Such end members close off the hollow ends, precluding each end from undesirably snagging other objects and further enhancing the appearance of the molding by concealing section components. The end members are generally of polymer material and are often differently shaped from each other. For example, the design of certain automobiles requires the front end of a front door molding to be tapered so that it does not catch on the rear edge the of front fender when the door is opened. The rear end of a rear fender molding is often shaped to match the shape of the rear wheel well. In any case, this arrangement results in a trim piece made up of the extruded strip and end members, which must be assembled and fastened to each other during manufacture, and having aesthetically undesirable parting lines between the strip and end members. Its hollow construction is considered advantageous however, for its economical use of materials and light-weight aspect.\nAnother known type of molding comprises a extruded piece of solid plastic, possibly having a decorative metal strip coextruded as part of its exterior face. The molding is usually fastened to an automobile using an adhesive such as double-sided sticky tape. The extruded strip may be cut to obtain a blunt end during manufacture, but if section components are to be concealed, end members are required as with the hollow-type molding described above. Further, if a piece with shaped ends is desired, end members must generally be used because cutting of the extruded plastic to obtain, for example, a tapered shape results in a cut plastic surface exposed as part of the decorative face of the molding, and this is generally unacceptable in the automotive industry.\nA trim piece of hollow construction, not requiring end members to be attached to its ends and which may be readily attached to a vehicle without the use of clips would be considered an improvement over these known types of moldings."}
{"text": "The following describes a conventional motor driver.\nFIG. 14 shows the structure of a conventional motor driver.\nIn the figure, a rotor 1010 has a magnetic field unit achieved by permanent magnets, and generates a rotation force according to mutual action of windings 1011, 1012, and 1013. A power supplier 1020 is composed of three upper power transistors and three lower power transistors, and supplies power to the windings 1011, 1012, and 1013. A position detector 1030 compares each of terminal voltages V1, V2, and V3, which are each from one terminal of the windings 1011, 1012, and 1013, respectively, with a common voltage Vc, and outputs a detection pulse signal FG in accordance with the result of comparison. A command unit 1040 outputs a speed command signal EC for controlling speed of the rotor 1010. In accordance with the signal EC, the switching controller 1050 outputs a PWM signal Wp for having the upper power transistors of the power supplier 1020 perform PWM operation. In accordance with the detection signal pulse FG and the PWM signal Wp, a distribution controller 1060 outputs upper distribution control signals N1, N2, and N3 and lower distribution control signals M1, M2, and M3 for controlling power distribution to the windings 1011, 1012, 1013. Accordingly, the power supplier 1020 supplies power to the windings 1011, 1012, and 1013, and has the motor perform PWM sensorless driving.\nA further structure is disclosed in Japanese Patent Application Publication No. 2001-346394 (p. 18, paragraph no. 0051), for having position detection performed stably in order to eliminate instability in accelerated turning operation caused by lag in position detection.\nA problem exists in these conventional motor drivers in that startup failure occurs easily. Startup failure occurs because the rotor 1010 is unstable in terms of position and rotates slowly at the initial startup, and therefore the back EMF (electromagnetic force) voltage that is induced in the windings 1011, 1012, and 1013 is low. Consequently, the position sensor 1030, which detects position based on the comparison results of the terminal voltage V1, V2, and V3 of the windings 1011, 1012, and 1013 with the common voltage Vc, detects erroneously.\nParticularly, when having the motor perform PWM driving, induced noise that is characteristic of PWM operation is superimposed on the terminal voltage in the detection phase. As a result, the probability of the position sensor 1030 detecting erroneously further increases due to this superimposed noise.\nA further conventional technique that attempts to deal with this problem is a method that fixes the position of the rotor in startup according to magnetic pull in a specific phase. However, this method is problematic because the motor driver takes longer to start up due to the additional time required to fix the position of the rotor.\nIn view of the stated problems, the object of the present invention is to provide a motor driver that enables stable PWM sensorless startup in PWM sensorless driving, taking into consideration the effects of noise that is characteristic of PWM operation."}
{"text": "1. Field of the Invention\nThe present invention relates to a sample detecting device, and more particularly, to a multi-angle and multi-channel detecting device.\n2. Description of Related Art\nLiquid crystal display panel detecting technique includes measurements of optical parameters such as chromaticity and brightness and plays an important role in the quality control of liquid crystal displays. In particular, with the increasing area of the display panels and the increasing processing speed, a fast and accurate detection has become more and more important. At present, conventional panel detecting devices can be classified into single-point Fourier optics panel measuring devices and multi-point beam-split panel measuring devices.\nFIG. 1 shows a technique disclosed in U.S. Pat. No. 6,804,001. The disclosed device is a combination of Fourier optics theory and a beam-split spectrum image spectrum measuring structure. Light rays emitted from a test object 2 pass through lens sets 6, 8 to reach a slit 16. After passing through the slit 16, the light rays are projected by a beam-split element 18 onto a two-dimensional photodiode array sensor 14. The element 18 is composed of several filters with different bandpass ranges. In addition, the device further includes a rotating mechanism 40 capable of rotating the slit 16 with different angles, so as to obtain chromaticity and brightness of a test point at different view angles. If information regarding the chromaticity and brightness of the entire test object is required, moving the probe or the sample in two dimension manner and combining the chromaticity and the brightness information at every point of the sample are necessary. However, this device or method spends a lot of time for performing the measurement, and can hardly be implemented for an on-line detection.\nFIG. 2 shows a multi-point beam-split panel measuring device disclosed in U.S. Pat. No. 5,751,420. The disclosed device uses a beam-split image spectrum measuring structure. A panel 42 is placed on a plane at a work distance from an image capturing lens. Light rays emitted from the panel are guided by the lens into a spectrum imaging device. After passing through the beam-split element inside the device, light signals of different wavelengths are projected onto different positions of a two-dimensional photodiode array sensor, so that the spectrum image corresponding to different positions in the field of view on the object side are obtained. Although this method can simultaneously obtain multiple channels of spectrum information to achieve a multi-point spectrum measuring effect, the measurement can only be carried out at one specified view angle at a time. Thus, a mechanism for providing an angular swing of the probe or of the panel is required when measurements at different view angles are performed. The rotating and positioning mechanism not only is complicated, but the motion itself is also time-consuming. As a result, although this method has been applied to on-line detection, it is not a fully satisfied solution.\nThe single-point measuring method is currently the most widely used panel inspection technique. To measure the entire panel, the probe or the panel has to be moved two-dimensionally and a lot of measuring time is required. Furthermore, in order to obtain optical parameters such as chromaticity and luminance at different view angles, the probe or the panel must be swung with a relative angle. However, the rotating and positioning mechanism can be quite complicated and the mechanical motion takes a lot of time. Now, there are still no rapid measuring devices that satisfy the VESA measuring standard, and thus an innovative measuring method is urgently required."}
{"text": "The present disclosure relates to vessels, and in particular to cup or bottles. More particularly, the present disclosure relates to a cup formed from polymeric materials."}
{"text": "1. Field\nThis disclosure relates to a red phosphor and a plasma display panel (PDP) including the red phosphor.\n2. Description of the Related Art\nPhosphors are materials that receive energy such as light or electrons from the exterior and emit light (from the visible light wavelength range) that is perceivable by the human eye. Phosphors may be used for such electronic devices as plasma display panels (PDP), field emission displays (FED), and light emitting diodes (LED), and they may dominate the light emission characteristics and color characteristics of the electronic devices.\nA PDP is a display device that displays an image by exciting phosphors with vacuum ultraviolet (VUV) rays generated by gas discharge in discharge cells.\nA phosphor layer of the PDP includes red, green, and blue phosphors. Such phosphors should exhibit satisfactory light emitting brightness, decay time characteristics, and lifetime."}
{"text": "Portable electronic devices, such as laptops, tablets, and smartphones, are increasingly used in settings in which sensitive information may be stored on the device. For example, a user may send and receive corporate email on their smartphone or may access, edit, and save internal corporate documents on a tablet device, such as via a remote network connection. In addition, it is becoming increasingly common for users to supply their own devices for use at work under policies commonly referred to as “bring your device” or BYOD policies. However, devices used under such BYOD policies present security risks for information technology (“IT”) departments, who traditionally provided corporate devices to users with particular security features enabled."}
{"text": "1. Field of the Invention\nThe present invention relates to floating fish baskets, and more particularly pertains to a new and improved floating fish basket for use by divers and spear fishermen to contain and protect from predators the day's catch of fish, lobster and other sea crustaceans. Spear fishing divers currently use plastic bottles, PVC plumbing pipe, inner tubes, foam tubes, bamboo trees, surf boards and a variety of other flotation devices attached to stringers or nylon nets to secure the day's catch of fish. A serious drawback of these types of stringer supports is that the catch is open to predators such as barracuda, sharks and eels. Nylon nets are disadvantageous because while although containing the catch, they are still susceptible to predator damage. Additionally, the above described devices are difficult to tow through the water. In order to overcome these problems, the present invention provides a lightweight, inexpensive, durable and easily transportable floating fish basket which utilizes a collapsible metal net to provide protection from predators.\n2. Description of the Prior Art\nVarious types of floating fish baskets are known in the prior art. A typical example of such a floating fish basket is to be found in U.S. Pat. No. 2,241,314, which issued to P. Mohler on May 6, 1941. This patent discloses a mesh bag suspended by lacing from a flotation ring. U.S. Pat. No. 2,739,410, which issued to F. Budnick on Mar. 27, 1956, discloses a collapsible container having an accordion fold bag suspended from a flotation ring. U.S. Pat. No. 2,790,266, which issued to L. Williamson on Apr. 30, 1957, discloses a floating live bait container which has a door secured by a spring extending between an inner surface of the door and the bottom of the container. U.S. Pat. No. 3,478,463, which issued to L. Ruter on Nov. 18, 1969, discloses a floating fish basket in which a wire mesh bag is suspended from a flotation ring and is adapted for quick release therefrom by removal of a plurality of retention clips. U.S. Pat. No. 3,524,278, which issued to H. Wolfe on Aug. 18, 1970, discloses a floating fish basket in which a wire mesh bag is suspended from a flotation ring. A pivotal door having a spring closure is provided. U.S. Pat. No. D. 257,377, which issued to D. Andrews on Oct. 14, 1980, discloses a floating live bait container with a generally hexagonal configuration. U.S. Pat. No. 4,570,374, which issued to G. Baxley on Feb. 18, 1986, discloses a floating fish basket suspended from a molded foam flotation frame. A plurality of inclined baffles within an opening of the frame retain fish therein.\nWhile the above mentioned devices are suited for their intended usage, none of these devices disclose a floating fish basket which utilizes a collapsible wire mesh bag suspended from a frame removably secured to a flotation ring and having a socket for supporting a diving flag. Additionally, none of the aforesaid devices disclose a floating fish basket having a mechanism for retaining a pivotal door closed and operative upon opening the door for partially axially collapsing a wire mesh bag. Inasmuch as the art is relatively crowded with respect to these various types of floating fish baskets, it can be appreciated that there is a continuing need for and interest in improvements to such floating fish baskets, and in this respect, the present invention addresses this need and interest."}
{"text": "Modern vehicle styling has removed typical dent protection, such as rubber bumpers, from the vehicle doors. Consequently, door dents and dings have increased in frequency, with the cost of the repair being shouldered solely by the vehicle operator. This not only increases the cost of vehicle ownership, but also the associated frustration level. Moreover, it is a problem incapable of redress by the vehicle owner without altering the appearance of the vehicle.\nFor instance, an example of a past proposal for a vehicle step that offers dent protection, but at the expense of aesthetics, is shown in U.S. Pat. No. 8,146,935. This arrangement is limited pragmatically because the mechanism for actuation is fully exposed to view, and thus presents an unattractive appearance that would be unacceptable to most vehicle owners (which are well known to be conscious of the appearance of their vehicle). Because the running board is also rotated into the protective position, it must be spaced a significant distance from the rocker, or else it will potentially engage and damage the vehicle on actuation (including if inadvertently is overextended). The step is also not in the form of a conventional running board, and thus could not be used with existing vehicle designs without significant modification.\nAccordingly, a need is identified for a new and improved active dent shield for a motor vehicle."}
{"text": "1. Field of the Invention\nThis invention relates to methods and systems for processing and responding to emergency medical inquiries. Specifically, this invention relates to the process or method for receiving and processing critical information regarding emergency medical calls. Also, this invention specifically relates to a system and apparatus for performing the steps of a process for receiving and processing critical emergency medical information.\nProviding adequate emergency medical care presents several critical challenges to medical care providers. These challenges include: the proximity to the care provider, the time required for help to arrive, the identification of the criticality of the emergency, the appropriate level of care provided, the variances in training of emergency medical dispatcher personnel, and limited nature of emergency care resources. This invention addresses these challenges by providing a consistent and proven system for: First, gathering necessary medical complaint information from emergency medical inquiry callers. Second, prioritizing the complaint to determine the criticality of the emergency. Third, providing emergency verbal instructions to individuals at the scene. Fourth, assisting dispatched responders to be prepared for each emergency situation. And, fifth, advising those on the way to provide care at the scene of specific problems or potential hazards. When used correctly this invention decreases the effective response time, while increasing the professionalism and control of emergency medical dispatchers; increases the accuracy and appropriateness of patient interrogation and well as the quality of gathered information; reduces the number of multiple unit and light-and-siren responses thereby reducing the risk of emergency medical vehicular collisions; improves patient care; reduces burn-out and stress of dispatchers by improving their quality of training and performance; decreases the risk of responder injury or mistake by providing responders with improved knowledge of the situation; and provides a means for continuously improving the quality of emergency medical dispatching and, as a result, emergency patient care.\n2. Description of Related Art\nIt is desirable to provide a systematic and standardized method for responding to emergency medical requests. Although in the related art some attempt has been made to address the problem of medical care assessment, the related art does not address the specific problems of emergency dispatcher response. Rather related art approaches known to the applicant describe the following. A process of helping patients assess their health, select appropriate health care, and guide such patients to an appropriate level and type of care. An automated medical history taking system and a technique wherein selected branch paths through a question repertory are provided. A method and apparatus for coordinating the actions of two or more medical teams, especially for instructional purposes. An expert system for providing suggested treatments for a patient with physical trauma. A medical payment system that incorporates computer technology in the storage, retrieval and processing of patient data and insurance claims. A knowledge base containing medical/pathological information on various diseases. A hospital computerized system for entering information pertinent to a patient's stay in the hospital. An expert computer system for processing medical claims. An interactive computerized apparatus and method for presenting medical information for diagnosis and study of disease. An automated and interactive positive motivation system to send a series of motivational messages and/or questions to a client to change or reinforce a specific behavioral problem. An artificial intelligent expert system. A rapid response health care communications system for providing rapid and reliable health services to patients located within or outside a health care facility.\nFor general background material, the reader is directed to U.S. Pat. Nos. 4,130,881, 4,237,344, 4,489,387, 4,839,822, 4,858,121, 4,945,476, 5,063,522, 5,065,315, 5,072,383, 5,253,164, 5,255,187, 5,471,382, and 5,596,994. Each of the above references is hereby incorporated by reference in its entirety for the material disclosed therein MICROFICHE APPENDIX. This specification includes a Microfiche Appendix which includes 1 page of microfiche with a total of 20 frames. The microfiche appendix includes computer source code of one preferred embodiment of the invention. In other embodiments of the invention, the inventive concept may be implemented in other computer code, in computer hardware, in other circuitry, in a combination of these, or otherwise. The Microfiche Appendix is hereby incorporated by reference in its entirety and is considered to be a part of the disclosure of this specification."}
{"text": "My invention relates generally to sound amplifier systems for musical instruments and more particularly to an optoelectronic amplifier in which the musical vibrations of the instrument are transmitted by reflected light rays to a photo-electric receiver device and amplified into audible tones.\nThe principle of transmitting sound by modulated and reflected light rays is old as evidenced by the early patent to Bell et al. U.S. Pat. No. 235,496. It is also well known to utilize a photoelectric cell to convert the modulated light rays into an electronic signal as shown in U.S. Pat. Nos. 3,065,352 and 3,733,953.\nConventional electronic amplifier systems for musical instruments possess several deficiencies which are not present in my system. Conventional amplification systems generally utilize a heavy magnetic core positioned within the instrument which not only adds additional weight to the instrument, but also tends to deaden the acoustics of the instrument. Further, conventional systems utilize a cord extending from the instrument to the amplifier units which is sometimes cumbersome for the musician.\nMy invention solves many of these problems by providing an amplification system for musical instruments wherein very little, if any, additional weight is added to the instruments so that the acoustical characteristics of the instrument is not changed by any significant degree. Further, my amplification system eliminates the need for a cord extending from the instrument to the amplifier or power locations."}
{"text": "The present invention relates to a wire electrode for use in a wire-cut electrical discharge machining process.\nWire electrodes for wire-cut electrical discharge machining are generally in the form of a wire of copper, brass, tungsten or the like and have a diameter in a range of from 0.05 to 0.3 mm. FIG. 1 of the accompanying drawings is illustrative of the manner in which electrical discharge machining is carried out with such a wire electrode. The wire electrode, designated at 1, is tensioned and fed along at a constant speed in the direction of the arrow A while being held in a confronting relation to a workpiece 2. Then, a machining solution 3 is applied in a direction coaxial with the wire electrode 1 while a pushed voltage is impressed between the wire electrode 1 and the workpiece 2. An electrical discharge is now repeatedly produced through the medium of the machining solution 3 across a small gap between the wire electrode 1 and the workpiece 2 to melt and scatter away an amount of material of the workpiece 2 by heat energy generated upon the electrical discharge. An XY crosstable (not shown) coupled to the workpiece 2 is numerically controlled to achieve desired relative movement between the wire electrode 1 and the workpiece 2 while keeping the electrode-to-workpiece gap constant at all times and ensuring continuous electrical discharge.\nBy repeating the electrical discharge and controlling the XY crosstable in the above manner, a groove 4 can be continuously cut in the workpiece 2 to machine the workpiece 2 to a desired contour. Such wire-cut electrical discharge machining has been widely used in blanking and cutting general dies, for example.\nThe speed of wire-cut machining is dependent on the degree of tension applied to the wire electrode 1, as shown in FIG. 2 where the abscissa indicates the tension T (g) and the ordinate the cutting speed F (mm/minute). FIG. 2 shows a characteristic curve which progressively rises as it goes to the right, the indication being that the cutting speed is higher as the tension is larger. It has been confirmed that as the tension is made larger, the wire electrode 1 is subjected to smaller vibrations and the electrode-to-workpiece gap can be controlled more uniformly for stabler electrical discharge repetitions, resulting in a higher cutting speed.\nWire electrodes of copper, brass or steel, for example, having conventional crystalline structures suffer a limitation on the tensile strength thereof, and it is not possible to achieve a higher cutting speed through an increase in tensile strength.\nWhen a conventional wire electrode 1 of copper, brass or steel is fed along upwardly or downwardly with respect to a workpiece during machining as shown in FIG. 3, portions of the wire electrode 1 are often scattered and deposited on an upper or lower end of a groove 4 cut in the workpiece 2. The deposited material 5 is mainly composed of copper or steel, and it has been observed that the material is deposited behind the wire electrode 1 as it cuts into the workpiece 2 as illustrated in FIGS. 3A. The deposit 5 on the machined surface tends to impair the dimensional accuracy of the cut groove 4. Such a deposited layer 5 has a thickness in the range of about 10 to 100 microns in areas where large machining energy is applied. As the machining energy is increased, the cut groove 4 is sometimes filled with the deposited material as shown in FIG. 4.\nThis undesirable phenomenon results in various drawbacks. The workpiece having been machined cannot be removed from the wire electrode. During machining, the machining solution 3 ejected coaxially with the wire electrode 1 does not enter the electrode-to-workpiece gap, causing a gaseous electrical discharge to lower the cutting speed and resulting in a danger of breaking the wire electrode 1. The deposit 5, mainly of copper, iron or the like, can only be removed with a dangerous chemical such as fuming nitric acid, a procedure which is tedious, time-consuming, and dangerous.\nTherefore, conventional wire electrodes have suffered from many drawbacks and have proven unsatisfactory."}
{"text": "In general, a cache is used to duplicate a certain part of a lower level memory, so that the duplicated part in the cache can be accessed by a higher level memory or a processor core in a short amount of time and thus to ensure continued pipeline processing of the processor core.\nCurrently, cache addressing is based on the following ways. First, an index part of an address is used to read out a tag from a tag memory. At the same time, the index and an offset part of the address are used to perform an addressing operation to read out contents from the cache. Further, the tag from the tag memory is compared with a tag part of the address. If the tag from the tag memory is the same as the tag part of the address, called a cache hit, the contents read out from the cache are valid. Otherwise, if the tag from the tag memory is not the same as the tag part of the address, called a cache miss, the contents read out from the cache are invalid. For a multi-way set associative cache, the above operations are performed in parallel on each set to detect which way has a cache hit. Contents read out from the set with the cache hit are valid. If all sets experience cache misses, contents read out from any set are invalid. After a cache miss, cache control logic fills the cache with contents from the lower level storage medium."}
{"text": "This invention relates to phosphoric acid esters and their salts, to paints and molding compositions containing them, to the use of these phosphoric acid esters and salts as dispersants, and to solids coated therewith.\nHigh mechanical forces are needed in order to incorporate solids into liquid media. To reduce these dispersing forces it is common to use dispersants which facilitate incorporation of the solids into the medium. These are generally surface active substances, also called tensides, having an anionic, cationic or non-ionic structures. These substances are either applied directly to the solid in small amounts, or they are added to the dispersing medium. The difficulty of dispersion is reduced by such a tenside. It is furthermore known that such solids can tend to reagglomerate after the dispersing process due to mutual forces of attraction, which nullifies the previously expended dispersing efforts and leads to problems.\nUnsatisfactory dispersion manifests itself in an increase in viscosity in liquid systems, loss of gloss and changes of color in paints and coatings, unsatisfactory coloring power in pigmented molding compositions, and reduction of mechanical strength in reinforced plastics.\nAlso, an important factor is the high viscosity of liquid systems having a high solids content. To operate economically, an attempt is made to keep the solids content of such solid-containing dispersions as high as possible in the dispersion, combined with the lowest possible viscosity of the system that is to be dispersed. This is of particular importance in liquid synthetic resin systems, such as unsaturated polyester resins, for example, in which a very low initial working viscosity is desired with the highest possible solids content.\nPhosphoric acid esters having various structures have been disclosed, for example, for dispersing or treating pigments. U.S. Pat. No. 4,456,485 describes acidic or neutral phosphoric acid esters of fatty alcohols and alkoxylated fatty alcohols as agents for treating specific blue pigments (alkali blue). Likewise, European Patent Application No. EP 256,427 describes the use of previously known phosphoric acid esters of alkoxylated fatty alcohols for preparing pigment dispersions suitable for aqueous applications. U.S. Pat. No. 4,720,514 describes pigment dispersions made using phosphoric acid esters of alkoxylates having various structures. German Patent Application No. DE 2,726,854 describes phosphoric acid esters of polyols used as dispersants in aqueous media or as detergent builders. These phosphoric acid esters are of a very complex nature and also contain free hydroxyl groups which originate from the epoxides that are used. They are evidently of a polyvalent nature with respect to the free phosphoric acid groups present in the molecules. U.S. Pat. No. 4,717,424 starts from carboxyl group containing phosphoric acid esters which in addition to the phosphoric acid group contain a free --COOH group in the molecule. These are not polymers. These products are used for stabilizing metal pigments against attack by water. U.S. Pat. No. 4,698,099 describes pigment dispersions containing phosphoric acid esters of monohydroxy-terminated polyesters as dispersants.\nIn may cases, good results can be obtained with the phosphoric acid esters described above. However, these compounds are not widely usable in modern binding resins and plastic systems, such as for example in high-solids paints, aqueous paints, systems which can be hardened by ultraviolet and electron beams, or LS-SMC or LP-SMC molding compositions (low-shrink or low-profile sheet molding compositions).\nIn a number of different systems, these compounds are not sufficiently compatible. This leads to inadequate dispersion due to insufficient interaction with the surrounding medium, and in many cases to precipitation phenomena and the associated dulling, spotting and increased viscosity."}
{"text": "Field of the Invention\nThe present invention relates to wireless communications, and more particularly, to a method and apparatus for configuring a frequency hopping pattern for a machine-type communication (MTC) user equipment (UE) in a wireless communication system.\nRelated Art\n3rd generation partnership project (3GPP) long-term evolution (LTE) is a technology for enabling high-speed packet communications. Many schemes have been proposed for the LTE objective including those that aim to reduce user and provider costs, improve service quality, and expand and improve coverage and system capacity. The 3GPP LTE requires reduced cost per bit, increased service availability, flexible use of a frequency band, a simple structure, an open interface, and adequate power consumption of a terminal as an upper-level requirement.\nIn the future versions of the LTE-A, it has been considered to configure low-cost/low-end (or, low-complexity) user equipments (UEs) focusing on the data communication, such as meter reading, water level measurement, use of security camera, vending machine inventory report, etc. For convenience, these UEs may be called machine type communication (MTC) UEs. Since MTC UEs have small amount of transmission data and have occasional uplink data transmission/downlink data reception, it is efficient to reduce the cost and battery consumption of the UE according to a low data rate. Specifically, the cost and battery consumption of the UE may be reduced by decreasing radio frequency (RF)/baseband complexity of the MTC UE significantly by making the operating frequency bandwidth of the MTC UE smaller.\nSome MTC UEs may be installed in the basements of residential buildings or locations shielded by foil-backed insulation, metalized windows or traditional thick-walled building construction. These MTC UEs may experience significantly greater penetration losses on the radio interface than normal LTE UEs. Thus, for these MTC UEs, coverage enhancement may be required. The MTC UEs in the extreme coverage scenario may have characteristics such as very low data rate, greater delay tolerance, and no mobility, and therefore, some messages/channels may not be required.\nFor efficient MTC transmission method, a new frame structure for MTC transmission has been discussed lately. Further, various enhancements of MTC transmission method based on the new frame structure for MTC have also been considered."}
{"text": "As computing devices become increasingly ubiquitous in everyday life, the use of public devices for sharing content is becoming increasingly prevalent. Different people may use these public devices throughout the day for a variety of tasks, including accessing or communicating information. While the availability of public devices is useful, it is not without its problems. One such problem is that access to a user's documents or programs can be slow, resulting in user frustration with these devices."}
{"text": "A gas turbine includes a compressor that compresses air, a plurality of combustors that combusts fuel in the air compressed by the compressor to generate combustion gas, and a turbine that is driven by the combustion gas. The fuel supplied to the combustors is often preheated before being supplied to the combustors to increase the combustion efficiency in the combustors.\nPatent Document 1 below discloses a gas turbine plant that preheats fuel supplied to a combustor. This gas turbine plant includes a gas turbine, a heat recovery boiler that generates steam by the heat of exhaust gas from this gas turbine, a steam turbine that is driven by steam generated by this heat recovery boiler, and a preheater that heats the fuel to be supplied to the combustor. A first steam line that supplies the steam generated by the heat recovery boiler as a fuel heat source of the preheater is connected to the preheater. Furthermore, a second steam line that introduces steam that heated the fuel into a transition piece of the combustor is connected to this preheater. In addition, a third steam line that introduces the steam that cooled this transition piece into the steam turbine is connected to the transition piece of the combustor. This gas turbine plant further includes an attemperation spray line that sprays water on the second steam line to cool the steam that flows through the second steam line, that is, the steam that heated the fuel."}
{"text": "1. Field of the Invention\nThe invention relates in general to a KVM switch, and more particularly to a KVM switch including an embedded multi-media system.\n2. Description of the Related Art\nAlong with development of technology, management policies of various computer systems are actively developed in today's market. For example, the well-known KVM switch has a port connected to more than two computer systems at a time, and an input/output port connected to a set of user-interface device, such as a keyboard, a mouse, a display and an audio device, e.g. a horn. The user can select a compute system via the KVM switch and control the computer system via a peripheral user-interface device of the KVM switch. However, how to improve the function and added value of the KVM switch is an essential subject of the relevant industrials."}
{"text": "1. Field of the Invention\nThe present invention relates to messaging systems and more particularly pertains to a new message display system for displaying a series of messages that allows a user to express their thoughts and personalities.\n2. Description of the Prior Art\nThe use of messaging systems is known in the prior art. U.S. Pat. No. 5,105,179 describes an electronic license display comprised of an LCD and capable of serving as a message center. U.S. Pat. No. 5,905,434 describes a communications system using an LCD unit capable of displaying various messages controlled remotely from the interior of a vehicle. Another type of messaging system is U.S. Pat. No. 5,579,008 describing an electronic license plate system.\nWhile these devices fulfill their respective, particular objectives and requirements, the need remains for a license plate communication system."}
{"text": "The present invention relates to modular sample processing apparatus kit that can be used to process samples that may contain one or more analytes of interest.\nMany different chemical, biochemical, and other reactions are sensitive to temperature variations. Examples of thermal processes in the area of genetic amplification include, but are not limited to, Polymerase Chain Reaction (PCR), Sanger sequencing, etc. The reactions may be enhanced or inhibited based on the temperatures of the materials involved. Although it may be possible to process samples individually and obtain accurate sample-to-sample results, individual processing can be time-consuming and expensive.\nA variety of sample processing devices and apparatus have been developed to process samples quickly and efficiently. Examples of some particularly useful sample processing devices may be found in, e.g., commonly-assigned U.S. Pat. No. 6,734,401 titled ENHANCED SAMPLE PROCESSING DEVICES SYSTEMS AND METHODS (Bedingham et al.).\nAlthough these sample processing devices are useful, in many instances they provide significantly larger numbers of process chambers (e.g., 96, 384 or more process chambers) than needed for a particular assay (in which as few as one process chamber and associated reagents may be needed). Furthermore, many of these sample processing devices contain reagents that are by far the most costly component of the sample processing device. As a result, using only a portion of the process chambers and the reagents contained in a sample processing device can be prohibitively expensive."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates generally to a multiple-input multiple-output (MIMO) technology for an antenna structure, and more particularly to an electronic device including the antenna structure.\n2. Description of the Related Art\nThe advance of mobile communication and antenna technologies are accompanied by an increase of the number of antennas included in electronic devices. However, as electronic devices are generally small-sized, the available spaces for mounting the antennas is limited.\nFurther, because antennas are normally mounted in the limited interior spaces of the electronic devices, the antennas are often arranged very close to each other. Consequently, these neighboring antennas may interfere with each other, for example, when signals of adjacent frequencies are transmitted and received, thereby lowering of the performance of the antennas.\nAdditionally, when a housing of an electronic device is formed of a conductive material (for example, metal), the radiation performance of the antenna may also be affected."}
{"text": "1. Field of the Invention\nThe present invention relates to a unique package assembly suitable for use in storing and shipping a plurality of product containers, and which is convertible to a shelf-ready product display by simply removing and discarding a section of the original assembly that is designed for easy removal therefrom.\n2. Description of the Related Art\nShipping packages, such as corrugated, cardboard and/or paperboard boxes, trays and similar assemblies are commonly used to contain a plurality of products or product containers for transport of such products from the manufacturer, to their eventual retail environment. Upon arrival in a retail environment it is desirable to store or display these products on shelves so that the goods and products are conveniently visible and accessible by a user. Thus, some package assemblies have been designed to be convertible from a shipping configuration to a display configuration (shelf ready packaging).\nFor example, U.S. Pat. No. 7,451,878 to Rochefort et al. discloses package assemblies that are convertible to displays. One assembly described therein is formed of an outer wrap portion and an internal divider/structural support portion. Openings are provided in the face panel of the container. Zipper pull tabs are also provided in the side face panels so that the face panel and the top half of the outer wrap portion may be removed to convert the package assembly into a dispensing container or a tray.\nAnother shipping/display assembly is disclosed in U.S. Pat. No. 5,657,872 to Leftwich et al., issued Aug. 19, 1997. The disclosed package assembly has a tray portion and a cover portion. The front side panel of the tray portion has a severable portion, which may be removed to allow products be extracted from the resulting opening. U.S. Pat. No. 6,168,027 to Esser, issued Jan. 2, 2001, discloses a shipping/display box having a tear-out segment spaced from all eight corners of the box. The tear-out segment can be removed at a store for display.\nSome packaging assemblies available from Delkor Systems Inc. are described in U.S. Pat. Nos. 5,887,717; 6,588,594; and 7,370,761 and include package assemblies having product support members with tapered sides and substantially open fronts (with or without a lip) which allows the support members to be shelf ready upon removal of a shrink wrap and/or top pad upon arrival of the package assembly in a retail environment.\nEach of the above referenced U.S. patents and any and all patents, applications and/or all other published documents mentioned anywhere in this application are incorporated herein by reference in their entirety.\nDespite the plethora of unique convertible package assembly/display systems that are available, there remains a need for improved systems that are even more dynamically useful, environmentally friendly and cost effective."}
{"text": "1. Field of the Invention\nThis invention relates generally to light emitting diode lighting fixtures and, more particularly, the invention relates to light emitting diode lighting fixtures with emergency back-up and scotopic enhancement.\n2. Description of the Prior Art\nAlthough receptive field sizes account for some of the differences in visual sensitivity across the retina, the sensitivity at a given retinal location can also vary. The human eye can process information over an enormous range of luminance (about twelve (12) log units). The visual system changes its sensitivity to light; a process called adaptation, so that it may detect the faintest signal on a dark night and yet not be overloaded by the high brightness of a summer beach scene. Adaptation involves four major processes:\n1. Changes in Pupil Size. The iris constricts and dilates in response to increased and decreased levels of retinal illumination. Iris constriction has a shorter latency and is faster (about 0.3 s) than dilation (about 1.5 s). There are wide variations in pupil sizes among individuals and for a particular individual at different times. Thus, for a given luminous stimulus, some uncertainty is associated with an individual's pupil size unless it is measured. In general, however, the range in pupil diameter for young people may be considered to be from two (2) mm for high levels to eight (8) mm for low levels of retinal illumination. This change in pupil size in response to retinal illumination can only account for a 1.2 log unit change in sensitivity to light. Older people tend to have smaller pupils under comparable conditions.\n2. Neural Adaptation. This is a fast (less than one (1 s) second) change in sensitivity produced by synaptic interactions in the visual system. Neural processes account for virtually all the transitory changes in sensitivity of the eye where cone photopigment bleaching has not yet taken place (discussed below)—in other words, at luminance values commonly encountered in electrically lighted environments, below about 600 cd/m2. Because neural adaptation is so fast and is operative at moderate light levels, the sensitivity of the visual system is typically well adjusted to the interior scene. Only under special circumstances in interiors, such as glancing out a window or directly at a bright light source before looking back at a task, will the capabilities of rapid neural adaptation be exceeded. Under these conditions, and in situations associated with exteriors, neural adaptation will not be completely able to handle the changes in luminance necessary for efficient visual function.\n3. Photochemical Adaptation. The retinal receptors (rods and cones) contain pigments which, upon absorbing light energy, change composition and release ions which provide, after processing, an electrical signal to the brain. There are believed to be four photopigments in the human eye, one in the rods, and one each in the three cone types. When light is absorbed, the pigment breaks down into an unstable aldehyde of vitamin A and a protein (opsin) and gives off energy that generates signals that are relayed to the brain and interpreted as light. In the dark, the pigment is regenerated and is again available to receive light. The sensitivity of the eye to light is largely a function of the percentage of unbleached pigment. Under conditions of steady brightness, the concentration of photopigment is in equilibrium; when the brightness is changed, pigment is either bleached or regenerated to reestablish equilibrium. Because the time required to accomplish the photochemical reactions is finite, changes in the sensitivity lag behind the stimulus changes. The cone system adapts much more rapidly than does the rod system; even after exposure to high levels of brightness, the cones will regain nearly complete sensitivity in ten (10 min) minutes–twelve (12 min) minutes, while the rods will require sixty (60 min) minutes (or longer) to fully dark-adapt.\n4. Transient Adaptation. Transient adaptation is a phenomenon associated with reduced visibility after viewing a higher or lower luminance than that of the task. If recovery from transient adaptation is fast (less than one (1 s) second), neural processes are causing the change. If recovery is slow (longer than one (1 s) second), some changes in the photopigments have taken place. Transient adaptation is usually insignificant in interiors, but can be a problem in brightly lighted interiors or exteriors where photopigment bleaching has taken place. The reduced visibility after entering a dark movie theater from the outside on a sunny day is an illustration of this latter effect.\nStudies suggest that the primary photoreceptor system for melatonin suppression is distinct from the rod and cone photoreceptors for vision. This action spectrum suggests that there is a novel retinaldehyde photopigment that mediates human circadian photoreception."}
{"text": "This invention relates to a new and useful insulation support means for metal buildings.\nIn the construction of metal buildings, rigid frames are spaced along the length of the building and have ground supported uprights with lateral rafter portions. Longitudinally extending Z-shaped purlins are secured on the top of these rigid frames for roof support and have sufficient vertical dimension that provides an area for insulation between them. Since there is considerable spacing between the rigid uprights and between the purlins, some type of support is necessary to hold the insulation under the roof covering. Prior insulation supports have been complex constructed frames or baskets that are hung in place. With the use of constructed frames, there is the disadvantage that they are expensive as well as difficult and time consuming to install. The basket-type supports have the disadvantage that they have to be specially manufactured to fit purlin spacing and are expensive."}
{"text": "Computers and computer software make a very important part of modern life and business. Specifically, for any business organization to stay competitive, it is imperative that it uses various computer based software to manage data, customers, employees, etc. Businesses use a number of different software for virtually every function of their daily operations, including payroll, customer relation management, accounting, inventory management, etc. Each of these various software generally use one or more underlying database to store data and a number of graphical user interface (GUI) based forms to interact with the computer user. Depending on the business model, there may also be one or more intermediate layers operating between such database and the forms presented to the end user. This is especially true in the case of software using the Internet or other similar platform to provide applications.\nFor example, various accounting systems, customer relation management systems, human resource management systems, etc., operate as multilayer applications that run in network environment, allowing multiple users to access the system from multiple locations. These applications use a storage layer or a data layer to store data that may be accessed by the various locations over the network. For example, a database server or a collection of database servers may be used to store data used by such an application. Each of the various locations use an upper layer, such as a presentation layer or an interaction layer, to present the data from users or other application, to interact with users and other applications to receive, provide and change data, etc.\nAn example of such an upper layer may be an Internet browser, a spreadsheet application, an active server pages (ASP) application, etc., where such an upper layer may query data for one or more users, receive data from one or more users, and communicate data to the storage layer for storage. Quite often it is possible that one or more of the multiple users may need to transfer data to the storage layer at nearly the same time. In such a situation, it is imperative that the consistency of data passed to the storage layer is guaranteed. For example, for an accounting system using such multiple layer architecture, one user of the upper layer may be posting a transaction involving a bank account, while at the same time, another user may be querying the storage layer to get information related to the same bank account. Such an attempt by multiple users to access the bank account form the storage layer may cause a database deadlock, either denying access to both users, or in a worse case, communicating erroneous information to one or more of such multiple users. Therefore, these is a need to ensure that users interacting with multiple layer applications via an upper layer are able to initiate and commit transactions between such upper layer and the storage layer without causing database deadlock."}
{"text": "1. Field of the Invention\nThe present disclosure relates to purified preparations of mammalian pluripotent stem cells, preferably human pluripotent stem cells, derived from corneal limbus tissue. In preferred embodiments, the pluripotent limbal stem cell lines are self-renewing and have the ability to differentiate into tissues derived from all three embryonic germ layers (endoderm, mesoderm and ectoderm). Methods for isolating pluripotent limbal stem cell lines and methods of their use are also disclosed.\n2. Description of Related Art\nIn early development, the ultimate source of all tissues in a mammalian embryo or fetus are stem cells. In the embryonic stage embryonic stem cells (ES cells) are totipotent and therefore capable of developing into all the cells of a complete organism. Cellular development occurs through several phases, including cellular proliferation, lineage commitment, and lineage progression, resulting in the formation of differentiated cell types. There are three main lineages that are derived from embryonic germ layers: ectoderm, mesoderm and endoderm. The ectoderm germ layer forms the epidermis of the skin, sense organs, nervous system, and spinal cord. The mesoderm germ layer forms smooth muscle, connective tissues, blood vessels, heart, blood cells and bone marrow, reproductive organ, excretory system, striated muscles, and skeletal muscles. Finally, the endoderm germ layer forms epithelial linings of respiratory and gastrointestinal tract, pharynx, esophagus, stomach, intestine, and other associated organs. ES cells are referred to as pluripotent stem cells because they can differentiate into almost all cell types in an adult organism.\nDuring the last decade there has been ongoing research on the isolation and use of ES cells and cell lines, which in addition to having the ability to develop into most of the specialized cells in the human body also have the capacity to divide and proliferate indefinitely in culture. ES cells are often referred to as pluripotent stem cells because they are not fixed in their developmental potentialities and can differentiate into many different cell types in vitro. Cultured ES cells that are highly pluripotent can form clumps of cells in suspension culture referred to as embryoid bodies. ES cells were isolated from humans relatively recently (Thomson et al., (1998) Science 282:1145-1147; Gearhart, (1998) Science 282:1061-62). In embryoid bodies derived from human ES cells it is possible to discern differentiated cells bearing markers of a wide variety of cell types.\nThe isolation of human ES cells offers the promise of a remarkable array of novel therapeutics. Biologic therapies derived from such cells through tissue regeneration and repairs, as well as through targeted delivery of genetic material, are expected to be effective in the treatment of a wide range of medical conditions. However, despite the enormous potential of these materials, serious ethical issues related to the use of human pluripotent stem cells derived from human embryos or from fetal tissue obtained from terminated pregnancies make stem cell research and treatments problematic. In addition, technical issues associated with the use of ES cells are problematic. Tissues or cells derived from ES cells are not ideal for use in medical treatments because generally the ES cells will not be derived from the patient who will ultimately be receiving the treatment. Use of autologous tissues is preferred for stem-cell-based therapies in order to avoid tissue rejection problems.\nI. Adult Stem Cells\nThe problems associated with human ES cells led many researchers to turn their attention to adult tissues as a possible source of undifferentiated stem cells with properties similar to those of ES cells or germ cells derived from fetal tissue. It was known that after birth and throughout adulthood a small number of specialized stem cells are retained in an organism for the replacement of cells and the regeneration of tissues. Indeed, adult stem cells (also referred to as “tissue-specific stem cells”) have been found in very small numbers in various tissues of the adult body, including bone marrow, (Weissman, (2000) Science 287:1442-1446), neural tissue (Gage, (2000) Science 287:1433-1438), gastrointestinal tissue (Potten, (1998) Phil. Trans. R. Soc. Lond. B. 353:821-830), epidermal tissue (Watt, (1997) Phil. Trans. R. Soc. Lond. B. 353:831), hepatic tissue (Alison and Sarraf, (1998) J. Hepatol. 29:678-683), and mesenchymal tissue. (Pittenger et al., (1999) Science 284:143-147).\nNevertheless, while some potential sources of adult stem cells have been identified, to date adult stem cells have not been found to be an adequate replacement for ES cells. First, adult stem cells can be difficult to isolate because they are usually present only in minute quantities in tissues that are often not easily accessible, and their numbers appear to decrease with age. Second, adult stem cells appear to be a less desirable source of cultured tissue than ES cells because they have a shorter life span and less capacity for self-renewal. Third, adult stem cells are believed to be tissue specific and not pluripotent, generally capable of giving rise only to new cells of a few types closely related to their tissue of origin.\nOne particularly notable difference between ES cells and adult stem cells is that ES cells in suspension culture are capable of forming aggregates of cells known as embryoid bodies. These embryoid bodies usually contain germ cells of all three lineages that differentiate into various lineage-committed tissues. Therefore, embryoid bodies can be useful in the preparation of different types of differentiated cells in culture. To date, no other isolated adult stem cell lines have been reported that are capable of forming structures similar to embryoid-like bodies in culture.\nRecently, however, it has been suggested that some adult stem cells have the capacity to be pluripotent. The most fully characterized are the hematopoetic stem cells known as bone marrow stromal cells or mesenchymal stem cells (Jiang et al., (2002) Nature 418:41-48). These were the first adult stem cells found to have pluripotent properties. Pluripotent adult stem cells have also been isolated from liver (U.S. Publ. No. 2003/0186439), mouse inner ear (Li and Heller, (2003) Nat. Med. 9:1293-1299), and amniotic fluid (Prusa et al., (2003) Hum. Reprod. 18:1489-1493). Pluripotent adult stem cells have also been recently described in many tissues such as skeletal muscle, brain, and intestinal epithelium (Howell et al., (2003) Ann. N.Y. Acad. Sci. 996:158-173). Still, while these papers report isolated or identified adult stem cells that are pluripotent, these “pluripotent” adult stem cells, unlike ES cells, differentiate into only a few lineages. In addition, none of the isolated adult stem cells reported to date appear to be capable of forming embryoid-like bodies in culture in a manner similar to ES cells.\nII. Corneoscleral Limbus\nSimilar to the other sources of adult stem cells referenced above, it is known that adult stem cells are present in the corneoscleral limbus of the eye. These cells participate in the dynamic equilibrium of the corneal surface and replace superficial epithelial cells that are shed and sloughed off during eye-blinking. Severe damage to the limbal stem cells from chemical or thermal burns, contact lenses, severe microbial infection, multiple surgical procedures, cryotherapy, or diseases such as Steven-Johnson syndrome or ocular cicatrical pemphigoid can lead to destruction of limbal stem cells and limbal stem cell deficiency which can lead to an abnormal corneal surface, photophobia, and reduced vision (Anderson et al., (2001) Br. J. Opthalmol. 85:567-575). This damage cannot be repaired without the re-introduction of a source of limbal stem cells (Tseng et al., (1998) Arch. Opthalmol. 116:431-41; Tsai et al., (2000) N. Engl. J. Med. 343:86-93; Henderson et al., (2001) Br. J. Opthalmol. 85:604-609). Thus limbal stem cells, with their high proliferative capacity, are clearly crucial for the maintenance of a viable ocular surface because they provide an unbroken supply of corneal epithelial cells necessary to maintain the equilibrium of the corneal surface (Tseng, (1996) Mol. Biol. Rep. 23:47-58).\nExperiments conducted in the 1980s first indicated the existence of limbal cells in the corneal epithelium (Schermer et al., (1986) J. Cell Biol. 103:49-62; Cotsarelis et al., (1989) Cell 57:201-209). Although it was later suggested that the transcription factor P-63 was a specific marker for human corneal stem cells, this marker is also expressed in other epithelial cells such as skin, and therefore is not specific to corneal stem cells. In addition, although P-63 expression has been shown to be principally limited to the basal limbal region in human corneas (Moore et al., (2002) DNA Cell Biol. 21:443-51), in mice expression of this transcription factor was maximal in paracentral cornea tissue rather than limbus (Moore et al., (2002) DNA Cell Biol. 21:443-451). Therefore, currently there is no known definitive stem cell marker for limbal epithelial stem cells.\nIt would be desirable to identify a source of adult stem cells that are capable of self-renewal in culture and that are pluripotent and ES cell-like in their ability to differentiate into cells of all three major lineages: ectoderm, mesoderm and endoderm. Further, it would be desirable to isolate and culture these adult stem cells, and to induce them to differentiate into various cell types."}
{"text": "Recently, mobile communication devices and small personal computers are mounted with a camera module capable of generating a digital image or a digital video corresponding to an outside light. A camera module includes an image sensor module changing an outside light to an image and a lens focusing the outside light to the image sensor module. A conventional camera module is disadvantageously fixed by the lens and the image sensor module, and a distance between the lens and the image sensor module cannot be adjusted to make it difficult to obtain a high quality image desired by a user.\nRecently, a distance between the lens of a camera module and the image sensor can be adjusted by a VCM (Voice Coil Motor) to obtain a very high quality image from the camera module. A conventional VCM is configured in such a manner that an image sensor module is arranged at a rear surface, or a rear side of a base, and an IR (Infrared) filter is arranged at the rear surface of the base.\nHowever, the IR filter recently suffers from a disadvantage in that, concomitant with increased resolution of an image sensor module, quality of an image or a video generated by the image sensor module greatly deteriorates in a case a foreign object such as dust is penetrated into the IR filter."}
{"text": "Needle devices are used throughout the medical industry for the injection and withdrawal of a wide variety of fluids and solutions into and from the human body. Because of the numerous potential hazards associated with the handling and manipulation of bodily fluids, and particularly blood, there are a number of known safety features that are frequently incorporated into various types of needle devices to protect the practitioner from accidental exposure to the needle.\nPrior safety needle devices include several disadvantages including having a retractable sheath requiring long stroke distances to activate the safety feature, multi-component retraction and locking elements, and conveying an undesirable significant force against a patient's skin during activation of the safety feature upon receiving an injection. Conventional retraction syringe assemblies often also do not incorporate reuse prevention features, and thus, the retraction mechanism of the syringe may be reset so the syringe barrel may be reused. The reuse of syringe assemblies without sterilization or sufficient sterilization is believed to facilitate the transfer of contagious diseases. Further, the retraction features of conventional syringes may also require the user to actively activate the retraction mechanism. Accordingly, the chance of human error in failure to activate or properly activate the retraction mechanism can lead to continued exposure of needles leading to needle stick injuries.\nSome known retracting sheath safety needle devices have been developed to include a single-use safety needle device assembly that obscures a substantial majority or an entirety of an injection needle from view before, during, and after an injection procedure. However, many injection procedures require that the practitioner know precisely the location and depth to which the needle is inserted in the patient's tissue to be sure that medication is delivered to an appropriate location. In addition, there exists a tendency for many users to falsely assume that they were “safe” from needle stick injuries, even in the non-locked initial state, due to the tip of the prior art retracting sheath safety needle devices being fully covered in an unlocked state.\nThus, there is a need in the art to provide a safety needle device having a passive activation mechanism that overcomes the deficiencies of the known retracting sheath safety needle devices and which allows for shorter stroke distance, ease of use, increased patient comfort, low part count, low part complexity, relatively compact design, and clear and unobstructed view of needle in an initial position."}
{"text": "Contemporary medical imaging depends largely upon the use of radioisotopes. One of the first clinically-employed radioisotopes was technetium (Tc). This element was first administered to a human subject in 1961 in the form of Na99mTcO4. Other radioisotopes, including halogens, such as 125I, 131I and 82Br, and isotopes of various metal radionuclides of lead, gallium, rhenium, arsenic and copper, have also been explored as potential imaging agents. Medical imaging is used in a variety of medical applications, including imaging of the brain, tumors, and components of the cardiovascular system.\nBlood flow imaging agents are currently the most important tool for determining heart function. T1-201, Tc-99-MIBI and Tc-99-tetrofosmin are in routine use for myocardial imaging at rest and after exercise. These agents are very useful but are not optimal. These tracers are Single Photon Imaging agents and their resolution is limited to the properties of SPECT imaging cameras and technology. However, fluorine-18 can be detected by Positron Emission Tomography imaging technology which has several advantages including higher resolution and corrections for the emitted radiation attenuation. In fact, the number of PET cameras and imaging centers are growing rapidly in response to the superior performance properties of fluorine-18.\nF-18 is one of the most useful positron emitting radionuclides currently being used in clinical nuclear medicine diagnosis. For example, 2-[F-18] FDG (2-[F-18]-fluoro-2-dexoy-D-glucose) is the radiopharmaceutical of choice for the diagnosis of several cancers and brain disorders. This radiopharmaceutical agent produces superior high-resolution images and quantitative regional uptake of tissues. The 110-min half-life of fluorine-18 allows production and distribution of 2-[F-18] FDG to nuclear medicine facilities near a cyclotron center. The relatively long physical half-life of fluorine-18 also permits PET studies of moderately slow physiological process. Decay of fluorine-18 is largely by positron emission (97%), and the emitted positron is of relatively low energy (maximum 0.635 MeV) and thus has a short mean range (2.39 nm in water). Fluorine-18 is readily available from both particle accelerators and nuclear reactors using a wide variety of nuclear reactions, and can be produced at specific activities approaching the theoretical limit of 1.171×109 Ci/mmol.\nIn addition to their superior medical imaging properties, fluorine atoms are a component of many pharmaceutical compounds. Fluorine can function as a substitute for a hydrogen atom in many biologically active molecules without substantially altering their properties, as done in the case of 2-deoxy-D-glucose.\nDespite the utility of F-18, there are only a very small number of methods to introduce F-18 into organic molecules. To date, the introduction of F-18 to a single bond was made via an exchange reaction on mesylate or triflate. Alternatively, F-18 could be introduced onto a C—C double bond or aromatic ring via an appropriate tin compound and [F-18]F2 or using anhydrous [F-18]fluoride on an electron withdrawing activated ring. The exchange reaction is carried out by treating the mesylate or triflate with a mixture of F-18, potassium carbonate, and crown ether, such as Kryptofix. 2-[F-18] FDG is the best example of that reaction. Other reactions using, [F-18]-F2, [F-18]XeF2, [F-18]DAST, [F-18]triethylammonium fluoride were also reported for specific radiolabeling. Radiofluorination of tributyltin-substituted double bonds and aromatic rings used [F-18]F2 as a reagent. However, the specific activity of these radiofluorinations is very low due to the cold F2 carrier. Radiofluorination of a nitro moiety on an activated aromatic ring with F-18 anhydrous fluoride was also reported. However, most fluorine-containing drugs are not activated with electron withdrawing groups, such as nitro, aldehyde, ketone, ester or others; therefore, this reaction is not applicable for a large number of compounds.\nThere is an urgent need for the development of new agents that can improve the diagnosis of heart disease by understanding the molecular behavior, physiology, anatomy, and function of the myocardium. However, many biologically-active molecules, drugs, receptor ligands, peptides, and proteins are not readily available for clinical nuclear medicine due to the limitations inherent in the methods used to install F-18. Therefore, the need exists for a new method for labeling a compound with F-18 which is amenable to a wide variety of organic substrates."}
{"text": "1. Field of the Invention\nThe present invention generally relates to optical pickup units and information recording and reproduction apparatuses, and more particularly to an optical pickup unit recording information on and reproducing information from information recording media and an information recording and reproduction apparatus including such an optical pickup unit.\n2. Description of the Related Art\nDisk-type optical information recording media widely used nowadays are compact disks (CDs) such as CD-ROMs, CD-Rs, and CD-RWs and digital versatile disks (DVDs) such as DVD-ROMs, DVD-Rs, and DVD-RWs. Some CDs have realized a recording density of 650 MB. The DVDs are larger in capacity than the CDs, but have yet to satisfy the demands of users in terms of capacity. In this context, a so-called multilayer disk, a disk formed of a plurality of recording layers instead of a single recording layer, has been developed for achieving higher recording density.\nIn such an information recording medium of multiple recording layers, the recording layers are required to be separated from each other by tens of micrometers (μm) or more so that information may be recorded on and reproduced from each recording layer independently. However, distances from an objective lens to the recording layers are different so that spherical aberration occurs in recording layers out of an optimum position. That is, in such recording layers, the position of the focus of marginal rays (rays entering the periphery of a lens away from the optical axis thereof) entering a lens 90 shown in FIG. 1 is deviated in the direction of the optical axis from the position of the focus of paraxial rays (rays entering the central part of a lens) entering the lens 90.\nFIG. 2 is a diagram for illustrating focusing of a light beam in the case of reproducing information from a multilayer disk by a single conventional optical pickup unit. In the case of reproducing information from a first recording layer 101a of a multilayer disk 101, the light beam is focused on the first recording layer 101a through an objective lens of the optical pickup unit at a position indicated by 100a as shown on the left side in FIG. 2. In the case of reproducing information from an nth recording layer 101n, which is farther away from a disk substrate surface 101s than the first recording layer 101a is, the light beam is focused on the nth recording layer 101n through the objective lens at a position indicated by 100b as shown on the right in FIG. 2.\nIf the optical pickup unit is optimized for focusing the light beam into a spot on the first recording layer 101a to reproduce the information therefrom, the optical pickup unit has no problem in reproducing the information from the first recording layer 101a. However, when the objective lens of the optical pickup unit gets closer to the disk substrate surface 101s of the multilayer disk 101 to reproduce the information from the nth recording layer 101n, and the light beam is focused into a spot on the nth recording layer 101n, spherical aberration occurs due to an interlayer thickness between the first and nth recording layers 101a and 101n. As a result, the spot formed on the nth recording layer 101n is larger in diameter than the spot formed on the first recording layer 101a as shown in FIG. 2.\nIn order to solve this problem, it is necessary to develop a new optical head control technology of performing aberration correction by measuring the amount of spherical aberration of a beam spot. Japanese Laid-Open Patent Application No. 2000-155979 discloses an aberration detection device as means for solving this problem.\nFIG. 3 is a diagram showing a configuration of the aberration detection device. As shown in FIG. 3, a light beam emitted from a light source 201 and reflected from an optical disk 206 is split by a half mirror 202 to be divided into a light beam passing through a specific region and a light beam passing through the other regions by a hologram 209. The light beam passing through the specific region is deflected by the hologram 209 to be received by a plurality of photodetectors 207 so that the photodetectors 207 obtain respective signals. The obtained signals are compared so that an aberration is detected. The detected aberration is transmitted via an aberration signal processing circuit 208 to an aberration correction device 204 so that the aberration correction device 204 can be driven in real time based on the aberration so as to reduce the aberration of the optical system. In FIG. 3, reference numerals 203 and 205 denote a collimator lens and an objective lens, respectively.\nThe assembly conditions of a light-receiving element in the optical pickup unit are extremely strict so that the light-receiving element is required to be provided with an accuracy of a few micrometers or less. This reduces yield rate, thereby affecting the cost to a considerable extent. Further, the aging and temperature characteristics of the light-receiving element are subject to change. Therefore, it is desirable that a light-receiving element pattern be simple. However, according to the technology disclosed in Japanese Laid-Open Patent Application No. 2000-155979, an incident light beam for aberration detection is split into a plurality of beams by a hologram so that a plurality of light-receiving elements detect a given one of the split beams. Such a configuration, where each split beam is focused into a small spot and a plurality of light-receiving elements detect a given one of the split beams, is complicated and may impair the stability of the optical pickup unit. This may reduce yield rate, thereby incurring an increase in cost.\nThe optical pickup unit is known as a device for recording information on and reproducing recorded information from an information recording medium. The ratio of the intensity of a light beam focused onto the information recording medium for recording to that for reproduction ranges from 5:1 to 15:1. Generally, the emission power of a light source switches between recording time and reproduction time substantially in accordance with the ratio. However, in some semiconductor lasers, noise characteristics worsen as an output lowers.\nThat is, part of a light beam emitted from a semiconductor laser returns to the semiconductor laser as a returning light. A different resonator other than the semiconductor laser is formed between the returning light and the information recording medium. As a result, the state of oscillation of the semiconductor laser becomes unstable so that the output of the semiconductor laser includes noise. Further, noise is also generated in the output of the semiconductor laser by the operation of the semiconductor laser or a variation in environmental temperature. When the light-emission power of the semiconductor laser is far above a threshold current level to reach tens of milliwatts (mW), the light emission of the semiconductor laser is stable, being hardly affected by disturbances. However, when the light-emission power of the semiconductor laser is around the threshold current level, the light-emission power is affected by disturbances including those caused by the above-described returning light, so that variations are caused in the light-emission power.\nAccordingly, in the case of recording information on (writing information to) the information recording medium or erasing information recorded thereon, the light emission of the semiconductor laser is hardly affected by disturbances since the light-emission power of the semiconductor laser is far above the threshold current level. However, in the case of reproducing (reading out) information from the information recording medium, the light-emission power of the semiconductor laser is normally set to a low level so that the semiconductor laser emits the laser beam with a power of a few milliwatts slightly over the threshold current level, for instance, a power of five milliwatts. In this case, therefore, the semiconductor laser is especially subject to the returning light to be unstable in emitting the laser beam. Accordingly, the signal is deteriorated by noise caused in the output of the semiconductor laser.\nJapanese Laid-Open Patent Application No. 9-27141 discloses an optical pickup unit to solve this problem. According to this optical pickup unit, an electro-optical element capable of controlling transmittance of light is provided in an optical path from a light source to a recording medium. The electro-optical element controls transmittance for light emitted from the light source to a low rate (value) at a time of reproducing information, and to a high rate (value) at a time of recording information.\nFurther, in an optical pickup unit, light emitted from a semiconductor laser is focused onto a surface of an information recording medium through a focus optical system, and a reflected light from the information recording medium is directed through a detection optical system to a light-receiving element. Generally, the detection signal (electric current signal) of the light-receiving element of the optical pickup unit is converted into a voltage signal by a current-voltage conversion amplifier housed in the optical pickup unit to be output to a signal processing circuit.\nIf the amplitude level of the signal output through the current-voltage conversion amplifier is too low, a problem is caused in information reproduction. Therefore, such a configuration is employed that the gain of the current-voltage conversion amplifier is switchable so that the output amplitude level thereof falls within a proper range. However, in recent years, it has been required for the optical pickup unit to accommodate a plurality of conditions so as to be suitable for a variety of types of information recording media and various recording and reproduction conditions. Accordingly, in order to perform information recording, reproduction, and erasure in compliance with a plurality of types of optical disks, such as a CD-RW (compact disk rewritable), a CD-DA (compact disk digital audio), a CD-ROM, and a CD-R (compact disk recordable), as performed by an optical information recording and reproduction apparatus disclosed in Japanese Laid-Open Patent Application No. 10-255301, it is necessary to make such a gain switching circuit shown in FIG. 4 suitable for more combination patterns of disk types and light powers. Consequently, the number of resistors employed in the gain switching circuit increases so that the speed of response of a signal is reduced.\nIn addition, recently, it has been proposed to increase the numerical aperture (NA) of an objective lens for focusing a light beam onto an information recording medium so as to achieve high recording density. This is because the diameter of a beam spot can be reduced by using an objective lens of a large NA. However, as the NA increases, the rate of increase of aberration also increases. That is, spherical aberration, whose primary cause is a substrate thickness error in the information recording medium, is proportional to the NA to the fourth power, and coma, whose primary cause is the inclination of the information recording medium to the optical axis, is proportional to the cube of the NA.\nJapanese Laid-Open Patent Application Nos. 10-20263, 9-128785, 11-259892, and 2000-155979 disclose techniques to correct and detect wavefront aberration (spherical aberration, coma, and astigmatism), which techniques are known as prior art for solving the above-described problems.\nHowever, optical pickup units having such configurations as disclosed in the above-described references need to have their costs reduced by component sharing and their assembly processes simplified."}
{"text": "1. Field of the Invention\nThis invention relates to heating systems including furnaces having fuel-fired burners, and, more particularly, to a control arrangement which provides a safety interlock between a stack damper control apparatus and fuel supply apparatus for such systems.\n2. Description of the Prior Art\nHeating systems employing furnaces having fuel-fired burners require a vent stack to conduct combustion products away from the burner. Automatically controlled stack dampers are generally used in the ventilation stacks to permit the stacks to be closed when the furnace is not operating to minimize heat losses when the furnace is not operating. However, for safe operation, it is necessary that the stack damper be open in advance of each operation of the burner. Accordingly, systems in which automatic dampers are used generally include a control arrangement which provides an interlock between the damper control mechanism and fuel supply apparatus of the system to assure that the damper is fully open before the burner operates and is closed after the completion of the operation of the burner.\nIn one known arrangement in which a primary burner control is conditional on and subsequent to the opening of a stack damper, a drive motor is energized in response to a request for heat to drive the damper to an open position, and limit switches complete the burner circuit and deenergize the drive motor when the damper reaches the fully open position. The drive motor is reenergized at the end of the heat run to move the damper to the closed position, a further switch deenergizing the motor when the damper reaches the closed position. Movement of the damper from the fully open position permits a limit switch to interrupt the burner circuit. A time lag is provided between the interruption of the burner circuit and the closing of the damper to allow volatiles to be purged from the furnace following operation of the burner.\nWhen operating properly, systems such as the type referred to above provide the desired interlock between the stack damper and the fuel supply apparatus. However, under certain failure conditions, such as the welding together of contacts of the limit switches, or, when cam-operated switches are used, the cams becoming loose and shifting out of place, the fuel supply valves may be energized while the vent stack is closed. Moreover, due to the interlock arrangement, the system may become locked out following a momentary power interruption, an undesirable condition.\nTherefore, it would be desirable to have a control arrangement for use in a heating system which provides a safety interlock between a stack damper control apparatus and fuel supply apparatus of the system which prevents operation of the fuel supply apparatus whenever the damper is in a position other than a fully open position, or in the event of an unsafe failure of the control apparatus, and which prevents lock out of the system following a momentary power interruption."}
{"text": "Dielectric insulation media in liquid or gaseous state are conventionally applied for the insulation of an electrical active part in a wide variety of electrical apparatuses, such as switchgears or transformers.\nIn medium or high voltage metal-encapsulated switchgears, for example, the electrical active part is arranged in a gas-tight housing, which defines an insulating space, said insulation space comprising an insulation gas usually with several bar pressure and separating the housing from the electrical active part without letting electrical current to pass through. Thus, metal-encapsulated switchgears allow for a much more space-saving construction than switchgears which are mounted outdoors and insulated solely by ambient air. For interrupting the current in a high voltage switchgear, the insulating gas further functions as an arc extinction gas.\nConventionally used insulation gases with high insulation and switching performance have some environmental impact when released to the atmosphere. So far, the high global warming potential (GWP) of these insulation gases has been coped with by strict gas leakage control in gas-insulated apparatuses and by very careful gas handling. Conventional environmental-friendly insulation gases like e.g. dry air or CO2 have a quite low insulation performance, so that gas pressure and/or insulation distances would have to be increased.\nFor the reasons mentioned above, efforts have been made in the past to replace these conventional insulation gases by suitable substitutes.\nFor example, WO 2008/073790 discloses a dielectric gaseous compound which—among other characteristics—has a boiling point in the range between about −20° C. to about −273° C., which is low, preferably non-ozone depleting and which has a GWP less than about 22,200. Specifically, WO 2008/073790 discloses a number of different compounds which do not fall within a generic chemical definition.\nFurther, U.S. Pat. No. 4,175,048 relates to a gaseous insulator comprising a compound selected from the group of perfluorocyclohexene and hexafluoroazomethane, and EP-A-0670294 discloses the use of perfluoropropane as a dielectric gas.\nEP-A-1933432 refers to trifluoroiodomethane (CF3I) and its use as an insulating gas in a gas-insulated switchgear. In this regard, the document mentions both the dielectric strength and the interrupting performance to be important requirements for an insulating gas. CF3I has according to EP-A-1933432 a GWP of 5 and is thus considered to cause relatively low environmental load. However, because of the relative high boiling point of CF3I (−22° C.) gas mixtures with CO2 are taught. Additionally, pure CF3I-gas has about the same insulation performance as conventional insulation media having a high insulation and switching performance, so that the proposed gas mixtures have around 80% of the specific insulation performance of a pure conventional insulation medium which would have to become compensated by increased filling pressure and/or larger insulation distance.\nTherefore there is an ongoing need for an insulation medium which causes even less environmental load than CF3I and does not require an increase of the gas pressure and/or the insulation distances above today usual values."}
{"text": "(1) Field of the Invention\nThe invention relates to computer systems. More specifically, the invention relates to error handling in an inheritance or nested hierarchies of arbitrary depth.\n(2) Related Art\nIn object-oriented programming, one object is frequently derived from another object with a third object being derived from the second object, and so forth. This is known as collaboration by inheritance. When an inheritance hierarchy is arbitrarily deep, the handling of errors at the various levels of inheritance becomes very complex. The complexity of error handling becomes exponentially proportional and, in the worst case, will be on the order of 2.sup.N where N is the depth of the hierarchy.\nWhen an error occurs in a prior art inheritance hierarchy, no clear protocol exists for handling the error. This leaves open the questions of whether the object handles the error by itself. Whether if at all, the parent is informed that an error has occurred. In other words, there is no uniform handling of errors in inheritance hierarchies. Rather, each hierarchy may handle errors in a unique way. This does not facilitate code reuse and ultimately results in longer and slower source code. Similarly, because there is no uniform way of handling errors, there is no uniform acts that can be expected when a severe error occurs, nor is there any procedure for notification of the application when an error can be bypassed or ignored. Moreover, from a programmer's perspective, it is often desirable to be able to determine at what level the error is generated and which objects are responsible for it. Prior art error handling on an ad hoc basis fails to address these issues.\nAccordingly, it would be desirable to have a system for uniformly handling errors in an arbitrarily deep inheritance hierarchy. It would be desirable to be able to identify an underlying source of a high level error to permit the underlying cause of the error to be treated."}
{"text": "This invention relates to crystalline, low melting .epsilon.-caprolactone polymers bearing basic amine functionalities which are linked to the ester chain ionically or covalently to induce catalyzed hydrolysis. The ester components can be derived from .epsilon.-caprolactone with or without small amounts of glycolide, and/or similar lactones. Such polymers with accelerated absorption profiles are especially adapted for use as transient coatings for absorbable multifilament surgical sutures.\nMultifilament surgical sutures such as Dexon.RTM. polyglycolide multifilament suture typically require a surface coating to improve the handling and knotting characteristics of the suture. Capitalizing on the desirable low melting temperature, crystallinity, and rheological properties of polycaprolactone as a coating material, several compositions based on this polymer were investigated as coatings for surgical sutures. Recognizing the fact that the .epsilon.-caprolactone homopolymer is essentially non-absorbable led to the development of copolymers of .epsilon.-caprolactone with variable amounts of more absorbable monomers to improve the coating absorbability. U.S. Pat. No. 4,624,256 discloses a suture coating copolymer of at least 90 percent .epsilon.-caprolactone and a biodegradable monomer and optionally a lubricating agent. Examples of monomers for biodegradable polymers disclosed include glycolic acid and glycolide, as well as well-known monomers typically used to prepare absorbable polymer fibers or coatings for multifilament sutures. U.S. Pat. No. 4,788,979 and U.S. Pat. No. 4,791,929 disclose a bioabsorbable coating of a copolymer of at least 50 percent .epsilon.-caprolactone and glycolide. Sutures coated with such polymers are reported to be less stiff than sutures coated with other materials and the physical properties of the coated suture are also reported to be acceptable. U.S. Pat. No. 4,994,074 discloses copolymers of a predominant amount of .epsilon.-caprolactone, the balance being glycolide and glycolic acid; the use of glycolic acid reportedly increasing the rate of absorption of the copolymer when used as a coating for multifilament surgical sutures.\nUnfortunately, the problem of adequate bioabsorbability of .epsilon.-caprolactone-based polymers without detrimental effects on their desirable properties as coatings still remains. Specifically, the use of sufficient amounts of glycolide to achieve sufficient absorbability of the copolymeric coating can compromise its crystallinity and melting characteristics, for it may become amorphous or liquid near room temperature. On the other hand, the strategy of using glycolic acid to achieve the reported results in coating absorbability does limit the ability to produce sufficiently long chain molecules to achieve optimum frictional properties, due to glycolic acid's known properties as both a ring-opening initiator and a chain terminator. Thus, a totally new approach to modifying the absorbability of polycaprolactone and its copolymers without affecting their desirable properties as suture coatings would be a more desirable goal."}
{"text": "Several solutions have become known to prevent the unlawful appropriation of the goods introduced and offered for sale in the different shops. Among these there are the safety systems in the case of which the products intended to be protected are equipped with passive aerials, so-called goods protecting labels, made of metal foil, fixed on a paper or plastic carrier, and in the shop, practically near the exits, there are receiver-transmitter partial units radiating signals at a determined frequency. When these are in operation, if an active goods protecting label gets into the protected area, the safety system alarms with the help of sound and light effects, and prevents the goods from being taken out of the shop unlawfully.\nSuch devices are described in patent specifications No. EP 704.591 and FR 2.727 549.\nHowever, in the case of products which are supplied for the shops in a cases, boxes or in some other packaging, there is the possibility of that the person intending to misappropriate the product tries the take it out of the selling place by leaving the packaging behind. Such products are, among other, tapes, CD-s or computer discs suitable for recording and storing audio and video signals.\nSafety devices have also been developed for the protection of video, audio and other data carrying devices. The essence of these is that the goods, with their original packaging, are put in a case that can be locked, practically made of transparent material, also accommodating the packaging. The partial unit which activates the safety equipment is placed in the internal space of the safety device, next to the goods intended to be protected.\nSuch safety devices equipped with a storing space are described among others in patent specifications No. WO 92/08026, U.S. Pat. No. 5,147,034 and WO 97/16615, while information relating to signalling units which activate the different safety equipment is contained among others in patent specification No. WO/96241 19.\nHowever, the defect of the known safety devices is that the separate signalling partial unit put in the internal space which accommodates the goods may be easily damaged in the course of putting the goods in the internal space and taking them out of there, and it may result in the failure of the alarm.\nIn the case of the signalling partial units connected to the protective case it is a problem that, especially in the case of the cases for the protection of the CD-s, the goods intended to be protected shades the signalling partial unit with radio-frequency, and so the partial unit is not able to emit signals of appropriate intensity and to turn on the safety device.\nIt must also be regarded as a disadvantage, that the manufacturing cost of the known signalling partial units put in the internal space of the protective case is relatively high, and it is still susceptible to breaking down.\nIt is also unfavourable that in the interest of reaching appropriate sensitivity the known signalling partial units are relatively spacious, and due to this in many cases they often cover information important for the customer on the packaging of the goods to be protected, and so their application may be a disturbing, purchase impeding factor."}
{"text": "This invention relates generally to the field of macroscopically manipulable nanoscale devices that permit information to be provided to or obtained form a nanoscale environment, and more particularly to the use of nanotubes attached to macroscale mounting members as nanoscale probes, fabricators and manipulators.\nThe development of mechanical, electrical, chemical and biological devices and systems that include or comprise nanoscale components, sometimes termed nanotechnology, has been slowed by the unavailability of or limitations inherent in devices that enable sensing, measuring, analyzing, and modifying objects with nanometer resolution and sensing, measuring, analyzing, moving, manipulating, fabricating and modifying objects with nanometer dimensions.\nOne class of devices that have found some use in nanotechnology applications are proximity probes of various types including those used in scanning tunneling microscopes (STM), atomic force microscopes (AFM) and magnetic force microscopes (MFM). While good progress has been made in controlling the position of the macroscopic probe to sub-angstrom accuracy and in designing sensitive detection schemes, the tip designs to date have a number of problems.\nOne such problem arises from changes in the properties of the tip as atoms move about on the tip, or as the tip acquires an atom or molecule from the object being imaged. Another difficulty with existing probe microscope tips is that they typically are pyramidal in shape, and that they are not able to penetrate into small xe2x80x9cholesxe2x80x9d on the object being imaged, and they may give false image information around sharp vertical discontinuities (e.g., steps) in the object being imaged, because the active portion of the xe2x80x9ctipxe2x80x9d may shift from the bottom atom to an atom on the tip\"\"s side. Moreover, conducting conventional probe microscope tips have never been successfully covered with an insulating material so that the only electrically-active element is the point of the tip itself.\nConventional probe microscope tips also are very rigid in comparison to many of the objects to be examined, and with xe2x80x9csoftxe2x80x9d samples (e.g., biomolecules like DNA) conventional AFM tips misrepresent the thickness of the object imaged, because that object is literally compressed by the action of the tip.\nThus, there is a need for macroscopically manipulable nanoscale devices for observing, fabricating or otherwise manipulating individual objects in a nanoscale environment that address the foregoing and other disadvantages of the prior art.\nThe present invention employs geometrically-regular molecular nanotubes (such as those made of carbon) to fabricate devices that enable interaction between macroscopic systems and individual objects having nanometer dimensions. These devices may comprise one or more individual nanotubes, and/or an assembly of nanotubes affixed to a suitable macroscopically manipulable mounting element whereby the device permits macroscale information to be provided to or obtained from a nanoscale environment.\nIndividual nanotubes or bundles of nanotubes can be recovered from a material (such as the carbon nanotube xe2x80x9cropesxe2x80x9d) grown by procedures described herein. Assemblies of nanotubes can be fabricated by physical manipulation of nanotube-containing material, or by self-assembly of groups of nanotubes, or by chemical, physical, or biological behavior of moieties attached to the ends or to the sides of the nanotubes or bundles of nanotubes. Individual nanotubes or assemblies of nanotubes can be grown to achieve specific characteristics by methods described herein.\nMore particularly, the devices of the present invention can comprise probes with tips comprising one or more molecular nanotubes. When attached to an appropriate motion transducer (piezoelectric, magnetic, etc.) the probe is capable of sensing, measuring, analyzing, and modifying objects with nanometer resolution and sensing, measuring, analyzing, moving, manipulating, and modifying objects with nanometer dimensions.\nA method for making such devices is disclosed, which includes the steps of (1) providing a nanotube-containing material; (2) preparing a nanotube assembly comprising at least one nanotube from the nanotube-containing material; and, (3) attaching the nanotube assembly to a macroscopically manipulably mounting element.\nThe nanoscale devices according to the present invention provide strong, reliably mounted probe tips and other nanoscale fabricators and manipulators, that are gentle, hard to damage, even upon xe2x80x9ccrashingxe2x80x9d into the working surface, that can be easily made electrically conductive, that can present a uniform diameter and precisely known atomic configuration, including precisely located derivitization with chemical moieties.\nThe devices of the present invention have a number of advantages over conventional microscopy probes (e.g. STM and AFM). A probe tip consisting of a single molecular nanotube or a few such tubes has the advantage that all its constituent atoms are covalently bonded in place and are unlikely to move, even under extreme stress, such as that occurring when the tip xe2x80x9ccrashesxe2x80x9d into the object being imaged. Moreover, the known, stable geometry of molecular nanotube tips allows one to more accurately interpret the data acquired by probe microscopes using such tips. In addition, molecular nanotubes are very compliant, buckling in a gently, predictable, and controllable fashion under forces that are small enough to avoid substantial deformation to delicate sample objects. Unlike currently used pyramidal probe tips, molecular nanotubes are very long with respect to their diameter, and can therefore reliably image the bottom areas of holes and trenches in the items being imaged.\nElectrically conducting nanotube tips can be coated with an insulating material to achieve localized electrical activity at the end of the probe element. This geometry facilitates probing of electrochemical and biological environments.\nMolecular nanoprobe elements have remarkably different chemical activity at their ends because the atomic configuration on the ends differs fundamentally from that of the sides. Consequently, one can selectively bond specific molecules to the tip end. This site specific bonding enables chemically-sensitive probe microscopy, and a form of surface modification in which some superficial atoms or molecules of the object being imaged react chemically with the probe tip or species attached or bonded to it. This delicate chemistry enables a form of surface modification that is not possible with conventional tips. This surface modification can serve as a direct manipulation technique for nanometer-scale fabrication, or as a method of lithography in which a xe2x80x9cresistxe2x80x9d is exposed by the chemical or electrochemical action of the tip."}
{"text": "The present invention is a means for temporarily connecting large diameter vacuum hose to ductwork, for the purpose of, but not limited to, establishing a vacuum within the ductwork during duct cleaning operations."}
{"text": "1. Field of the Invention\nThis invention relates generally to analog-to-digital converters. More specifically, this application relates to reduced latency in an analog-to-digital converter. The disclosure is associated with high-speed pipeline analog-to-digital converters (ADC's), and may also be extended to other applications and ADC architectures. For example, the techniques described herein may be applied to algorithmic ADCs, also known as cyclic ADCs, to delta-sigma ADCs, to SAR ADCs, etc., and generally to any ADC that may require one or more digital-to-analog converters.\n2. Description of the Related Art\nPipeline analog-to-digital converters (ADCs) are sub-ranging data converters in which a signal is quantized coarsely in several steps and the results of the different steps are then combined to achieve a high level of quantization accuracy. Pipeline ADCs are popular because they may operate at very high speeds (hundreds of MHz, depending on the technology employed) while achieving relatively large dynamic range.\nReferring to FIG. 1A, a typical pipeline ADC 10 includes a plurality of stages 12, 14, . . . , L. In the first stage 12, the input Vin is converted using a flash ADC (not shown) and is combined with results from the subsequent stages 14, 16, . . . , L to form an output. As will be discussed in greater detail below, the error in each stage is determined by converting the output of that stage using a digital-to-analog converter. The difference between the input to the stage and the error is the “residue.” The residue for each stage is amplified and fed to the next stage and converted in the same fashion in the next stage. The output of the last stage L is provided to a backend ADC 16 which resolves the last bits. All outputs may be provided for time alignment and digital error correction 20.\nAs shown in FIG. 1B, a typical stage 100 in a pipeline ADC includes an input signal sampling network 108, an N-bit coarse ADC 102 with its own sampling network, an N-bit digital-to-analog converter 104 (DAC) and an amplifier 106. The sampling network 108 and flash ADC 102 sample the input signal 101 at the same time. The DAC 104 converts the quantized signal back into analog form and this signal is subtracted 110 from the main input signal 101. The residue 105 resulting from this operation is amplified using amplifier 106 in order to occupy, typically and in the absence of errors, a portion of the range of the following stage, for example, half. Ideally, the residue then consists only of quantization noise.\nThree factors may limit the performance and speed of operation of a pipeline ADC: errors occurring within the DAC 104 known as element mismatches; errors occurring within the amplifier 106 due to gain and nonlinearity; and excessive delay through the flash ADC 102 and DAC 104 signal paths. All three result in degradation of the ADC linearity and signal-to-noise ratio (SNR).\nFIG. 2 illustrates pipeline stage using a prior technique for addressing DAC 104 errors. The stage 200 includes an input signal sampling network 208, an N-bit coarse ADC 202 with its own sampling network, an M-bit digital-to-analog converter 104 (DAC) (where M>N) and an amplifier 206. Processing elements 232, 234 are added between the coarse ADC 202 and the DAC 204 and hence increase the delay in that path. This delay is critical in high speed operation, namely operation at rates of 250 MSPS and higher. An application of the technique of\nFIG. 2 has been proposed by others to address the problem of element mismatches in flash ADCs, by permuting reference thresholds presented to the comparators to improve the linearity of an ADC.\nDAC and amplifier errors as described above are sometimes estimated and canceled or corrected using Dynamic Element Matching (DEM) and Harmonic Distortion Correction (HDC) techniques.\nDEM takes thermometer-coded outputs of a coarse ADC and permutes them before they are provided to the connected DAC elements. The permutation matrix is such that every ADC output can reach every DAC input. The method of permutation sometimes randomizes the DAC errors, thereby creating a white spectrum, or shapes the errors such that the energy of an error signal occupies a region of frequencies outside the band of interest. In FIG. 2, the pipeline stage includes a DEM block 232 between the coarse ADC 202 and the DAC 204 and the proper placement relative to the stage digital outputs. The DEM block 232 can be implemented efficiently through the use of transmission gates. However, it introduces a finite delay which is non-negligible at high clock rates.\nIn a pipeline ADC, DEM is used with additional digital processing that estimates the DAC error signal and effectively removes it from the output. If this were not done, the DAC noise would decrease the SNR. This estimation and removal of the DAC errors is referred to in the literature as DAC noise cancellation (DNC).\nContinuing with FIG. 2, a signal Σt is added 234 to the output of the coarse ADC 202. This signal consists of the sum of several random, independent sequences that are used in the estimation of the amplifier 206 errors, for example, gain and nonlinearity. The number of sequences depends on the order of nonlinearity that needs to be estimated: one sequence for linear gain error, three sequences for harmonic error, and so on.\nAlso shown in FIG. 2 is backend ADC 212, whose output is provided to amplifier 220 and HDC module 222, whose output is summed 230 with the output of DNC module 226.\nIn an HDC technique, the output of the residue amplifier 212 contains terms in the quantization noise of the coarse ADC 202, the random sequences and their interaction through the amplifier nonlinear characteristic. If the highest significant order of nonlinearity in the amplifier is 3, the output of the residue amplifier 212 contains one term proportional to a3 (Σt)3 where a3 is the third order nonlinearity coefficient and Σt=t1 +t2 +t', three random sequences that can each take on values +A or −A where A is a constant quantity. Therefore, Σt is a four level signal that can take on values −3A, −A, +A, +3A. Since the product of random independent sequences is also a random and independent sequence, multiplying the (digitized) residue amplifier output by (t1, t2 , t3 ) randomizes all terms except the one in a3 (Σt)3 which can be extracted with a lowpass filter.\nA consequence of adding the random sequences to the output of the coarse ADC 202 is that the word length increases and the DAC 204 size and complexity increase accordingly. That is why the DAC 204 resolution M is greater than the coarse flash ADC 202 resolution N. In a typical implementation, M=N+3.\nFIG. 3 is an example of a prior art implementation of a coarse ADC with N=2 bits resolution. This implementation is often called a flash ADC. Four comparators 302.n compare the input voltage on line 304 to four threshold voltages (THR1, THR2, THR3, THR4) respectively. In some implementations, the threshold voltages may be associated with a resistor ladder 306. Other voltage divider techniques may be used as well. If the input voltage is greater than THR1, then comparator 302.1 outputs a logical 1, otherwise it outputs a logical zero. Similarly, comparator 302.2 compares the input voltage with THR2 and so on. The output of the coarse ADC 206 is a digital word formed from the outputs of all the comparators. This word is often denominated a “thermometer code.” The number of logical ones contained in the thermometer code is the digital representation of the analog input voltage on line 304.\nOne example circuit embodying the comparator function is shown in FIG. 4. Only the circuitry to compare the input against THR1 (corresponding to 302.1) is shown for clarity of explanation. The circuit 302.1 includes eight switches 408, 410, 412, 414, 416, 418, 420, 422, two capacitors 402, 404, and a comparator 406. The switches are driven by periodic clocks denoted phase 1, phase 2. When phase 1 is true, phase 2 is false and vice versa. When phase 1 is true, the upper capacitor 402 is charged to the voltage THR1, whereas the lower capacitor 404 is charged to the input voltage. When phase 2 is true, the capacitors are connected to comparator 406. Many alternative circuit realizations of the comparator 302.1 exist in the relevant art, which accomplish the following function:Output 1=1 if (input−THR1)>0,0 otherwise\nThat is, the signal output 406 is a logical 1 if the input voltage is greater than the threshold voltage THR1, a logical zero otherwise."}
{"text": "Exemplary embodiments pertain to the art of electrical power generation and supply systems.\nSome systems, such as unmanned aerial vehicles (UAV's) or the like often utilize electrical power for propulsion and operation of onboard systems. Some such systems, such as medium-sized UAV's that require power levels in the range of about 1 KW to 30 KW, have relatively short mission times because the energy density of batteries is far too low to effectively work in this power range, and conventional internal combustion engines and jet engines are very inefficient at these low power levels. One option that has been developed is a tethered UAV system in which the UAV is connected to a power source on the ground by a tether. Use of a tethered UAV allows for an increase in mission duration time, but reduces an operating height and distance in which the UAV may operate, due to the constraint of the tether.\nFurther, power generation and supply systems in, for example, UAV's or the like, must often operate in dusty environments, which may effect operability of the power generation and supply system."}
{"text": "Pulsed RLG (radar level gauging) is becoming a more widely spread method for industrial level gauging, due to its simple and cost effective microwave components. The pulse in a pulsed RLG can be modulated by a carrier frequency (typically 6 or 24 GHz), or be an unmodulated DC pulse. In the latter case, it is common to use some kind of transmission line (coaxial line, twin line, etc.), sometimes referred to as a probe, is usually used to guide the electro-magnetic signal through the material in the tank where it is reflected by one or more interface surfaces (such as air/liquid) between different parts of the tank content. In the former case a transmission line or wave guide can also be used, but generally an antenna is used to form a vertical radar beam which is reflected at possible interface surfaces.\nPulsed radars typically apply different types of directional coupling. An example of a directional coupler including high speed sampling capabilities is described in U.S. Pat. No. 5,517,198. Directional coupling divides the available power between transmission line and receiving line, thus introducing significant attenuation of both transmitted and received signals, degrading the sensitivity of the system. This is in particular a problem for systems using a DC pulse, since the choice of directional coupler is limited by the extreme bandwidth of such a pulse, which includes also large wavelengths.\nFor this and other reasons, pulsed systems therefore typically have lower sensitivity compared to frequency modulated continuous wave (FMCW) radar. The sensitivity (ability to detect weak reflections) is an important virtue for any RLG as a high sensitivity may enable the use of a smaller antenna or a longer transmission line, all other parameters held constant.\nOne way to provide a directional coupling without sensitivity losses is to use a circulator, such as a ferrite circulator. However, such solutions are expensive, and their performance is typically temperature dependent, making them unsuitable for use in radar level gauges."}
{"text": "1. Field of the Invention\nThe present invention relates generally to an automotive transaxle (such as that suitable for use in an \"FF\" type, viz., a front engine front wheel drive, automotive vehicle) and more specifically to a final drive arrangement for such a transaxle which minimizes the engine shaking effect produced when torque is applied to the driven wheels of the vehicle.\n2. Description of the Prior Art\nIn previously a proposed \"FF-transaxle\" arrangement 1 such as shown in FIGS. 1 and 2 of the drawings, a transversely mounted engine and transmission unit 2 has been mounted on the vehicle chassis vehicle through elastomeric mountings 3 and buffer rod 4. However, this arrangement has suffered from the notable drawback that upon application of torque to the wheels 5 via a final drive unit 6 and drive shafts 7, (such as during acceleration, braking etc.) a reaction is produced which is transmitted back through the final drive unit 6 to the transmission and engine unit 2 causing same to displace the unit against the bias of the engine mountings 3.\nTo overcome this problem the hardness of the elastomer used in the engine mountings 3 has had to be selected somewhere between that suited for vibration damping and that suited for pitching or other relatively large amplitude movements of the engine in order to maintain a suitable degree of engine vibration damping while preventing the engine and transmission from being displaced relative to the vehicle chassis. This compromise however, has not enabled the desired amount of engine vibration damping which can only be achieved through the use of elastomers softer than that possible with the prior art arrangement and has not enabled the achievement of the desired riding characteristics of the vehicle."}
{"text": "Heretofore, self balancing hoists have been primarily pneumatic. Pneumatic or air balanced hoists are useful in many industrial applications in which a relatively heavy load such as a work piece or tool is suspended at a certain height and manually repositioned by the operator. The tool or work piece is raised or lowered by a pneumatic piston. An air pressure regulator controls the air pressure to bias the piston and thereby raise or lower the load. These self balancing pneumatic hoists have a means or regulator for adjusting the air pressure to balance the load in a suspended position. Once balanced, the load may be easily manipulated by an operator by applying a small upset force in the desired direction of travel."}
{"text": "A non-aqueous electrolyte secondary battery such as a lithium-ion secondary battery has high energy density. Recently, therefore, the non-aqueous electrolyte secondary battery is widely used as batteries for use in apparatuses such as a personal computer, a mobile phone, and a portable information terminal.\nIn order to enhance properties (such as safety) of the non-aqueous electrolyte secondary battery, various attempts have been made to improve a separator which is provided between a positive electrode and a negative electrode. In particular, a porous film made of polyolefin is excellent in electrical insulating property and exhibits good ion permeability. Therefore, such a porous film made of polyolefin is widely used as a separator for a non-aqueous electrolyte secondary battery, and various proposals relating to the separator have been made.\nFor example, Patent Literature 1 proposes a non-aqueous electrolyte battery separator that is a multilayer porous membrane (i) in which a porous layer that contains an inorganic filler or a resin having a melting point and/or a glass transition temperature of 180° C. or higher and has a thickness of 0.2 μm or more and 100 μm or less is provided on at least one surface of a polyolefin resin porous membrane and (ii) which has air permeability of 1 second to 650 seconds/100 cc.\nPatent Literature 2 proposes a non-aqueous electrolyte battery separator which is a separator including a polyolefin layer and a heat-resistant insulating layer that (i) is provided on one surface or each of both surfaces of the polyolefin layer, (ii) contains a heat-resistant resin and oxidation-resistant ceramic particles, and (iii) contains the oxidation-resistant ceramic particles at a ratio of 60% to 90%."}
{"text": "1. Field of Invention\nThis invention relates to a display tray as for a commercial type of a food display cabinet the tray being self formed of a blank.\n2. Description of the Previous Art\nThe tray structure herein is particularly adapted for use in a food case or cabinet such as in refrigerated cabinets as in grocery stores, for an orderly separation and display of goods which are commonly visible through see-through doors. It has been noted that there appears to be an absence of such trays in general use."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus for measuring a particle size and a method for measuring the particle size.\n2. Description of the Related Art\nAs conventional methods for measuring the particle size, there are known an electric detection method of electrically detecting a volume of the particle which is suspended in an electrolytic solution (see U.S. Pat. No. 3,757,213, for example) and an image analysis method of optically capturing an image of the particle and analyzing the size of the particle from the captured image (see Japanese Unexamined Patent Publication No. HEI 8(1996)-136439, for example).\nAs another method of measuring the particle size, there is known a method of estimating the particle size from light scattering characteristics of the particle. However, it is not easy to find out a light scattering characteristic such as shows a good and simple correspondence to the particle size, and therefore, a problem lies in that complicated calculation is required for obtaining the particle size from the scattered light."}
{"text": "1. Field of the Invention\nEmbodiments of the present disclosure relate to cooling fans, and more particularly to a fan control circuit.\n2. Description of Related Art\nAn electronic device usually requires cooling of the power supply or one or more components of the device. Traditionally, a fan may be used to cool the device where continuous rotation of the fan dissipates unwanted heat away from the device. A fan usually continuously rotates whenever the power supply of the equipment body is turned on, which may lead to undesired results. Unwanted noise from the fan rotation may be present in an environment where quietness is required. Additionally, continuous rotation of the fan will inevitably lead to accelerated deterioration in the equipment.\nWhat is desired, therefore, is to provide a fan control circuit which controls the rotational speed of the fan to efficiently dissipate heat."}
{"text": "1. Technical Field\nThe present invention relates to data processing systems and, in particular, to personal computing devices. Still more particularly, the present invention provides a method, apparatus, and program for improved user input on personal computing devices.\n2. Description of Related Art\nMobile computing devices, such as mobile phones, hand-held computers, and personal digital assistants, are now ubiquitous in today's society. These devices now typically have processors, memory, and storage that provide users with computing power nearly on par with that of desktop computers. Many users find, however, that the primary drawback of such devices is that of the input mechanisms. Manual entry of any significant amount of data into a mobile phone or personal digital assistant (PDA) is certainly a non-trivial task.\nText input is typically accomplished by one of three input mechanisms. A first text input mechanism is a hardware key mechanism. Mobile phones usually have a numeric keypad with which users may enter textual characters using combinations of button presses. This method of input is slow and cumbersome, since the task of entering a long string of text requires a considerable number of button presses. Other mobile computing devices include a full keyboard. However, with the small size of the device itself, the size and spacing of the keys renders the keyboard difficult to manipulate.\nMany mobile computing devices include a touch screen display device and a stylus that may be used to manipulate graphical elements on the display. Some attempts have been made to allow users to use handwritten characters on the touch screen display with varying degrees of success. Graffiti® is a handwriting recognition mechanism that recognizes simplified gestures for textual characters. Still, the user must write individual characters and the recognition mechanism has limited success.\nSome mobile computing devices provide a “virtual” keyboard that presents a miniature version of a keyboard on a touch screen. The user may then use keyboard layout to tap the virtual keys to enter text. However, tapping the keys remains a tedious hunt-and-tap operation using a stylus, which is not as natural as two-handed typing on a full size keyboard. Furthermore, the user must successfully depress each key individually by tapping with the stylus. This operation is tedious and difficult for many users."}
{"text": "The prior art discloses systems where RFID transponders are imbedded in network plugs and used to identify a network plug when it is inserted into a receptacle of an item of networking equipment. One drawback of such systems is that the networking equipment must be equipped with a system for powering the RFID transponders and sending identification information to an asset management system. Legacy systems typically are not equipped to power and communicate with the RFID transponders. As result, such legacy systems must typically be replaced.\nWhat is needed therefore, and an object of the present specification, is a system which can provide identifiable interconnections without replacing the networking equipment."}
{"text": "1. Field of the Invention\nThe present invention relates to software distribution, particularly electronic software distribution.\n2. State of the Art\nElectronic software distribution, or ESD, is technology that enables software publishers and resellers to use the Internet as a direct channel to customers. ESD provides a fast, efficient alternative to the distribution of shrink-wrapped software products in retail stores. ESD follows one of two models, either Buy-Before-You-Try (Buy/Try) or Try-Before-You-Buy (Try/Buy). Try/Buy is technologically more complex but is favored by buyers and is advantageous for sellers as well.\nA software tool for Try/Buy \"self-wrapping\" of software programs is sold by Preview Software, the present assignee- Using this tool, a software program is \"injected\" with a code module that performs trial protection and enables electronic purchase of the software program. The injected software program and related files are then packed as a program archive, or setup file set, using a conventional installation tool such as InstallShield.TM. installation software. A secondary setup executable is added outside the program archive and functions to install distribution information and instructions.\nThe foregoing process may be more clearly understood with reference to FIG. 1. A program executable has added to it trial and commerce information, resulting in a modified executable. The modified executable and other files are processed by an installer tool to produce a setup file set, e.g., setup.exe, that is a ready-to-install Try/Buy application.\nVarious other competing and incompatible Try/Buy ESD technologies have emerged, offered by such technology providers as TestDrive Software, Release Software and TechWave, Inc. ESD merchants have aligned themselves with different ones of these technology providers. As a result, software publishers have faced either choosing one technology provider and ignoring merchants aligned with other technology providers, or producing multiple different versions of a software title for different channels, resulting in considerable expense and duplication of effort."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a content distribution system, and more particularly, to system and method to remotely manage and audit set top box (referred to as “STB”) resources.\n2. Background of the Invention\nMany dwellings including single homes, apartments, condominiums, town house and lofts, have installed a media distribution system to improve the reception of audio and video content. Often, these distribution systems take the form of either a cable television system or a satellite television system. These systems receive a signal from a content supplier, such as a cable television company or a satellite television company. Generally, content suppliers offer a variety of services including different levels of subscriptions, pay channels and pay per view programs. Additionally, newer STBs, like Tivo, Replay TV and Ultimate TV, include computing resources including processors, hard disk drives, controllers, and operating systems.\nSome customers, in an effort to access unauthorized channels or services, attempt to modify STBs. For example, there have been problems with users opening their STBs and installing larger disk drives, and installing memory cards that provide the codes necessary to access unauthorized channels and services. Other forms of unauthorized use of STBs and the resources associated with STBs have also occurred."}
{"text": "In the marketing of certain types of products such as cosmetics, it can be difficult for a customer to select which of a wide variety of available items is most appropriate. For example, the particular facial makeup base which is appropriate may depend upon the oilyness of a person's skin and the particular color of skin, and yet many persons may not know what terms best describe the oilyness or color of his or her skin, or how to select an appropriate makeup base even if they could describe it. Even if a manual were provided, many customers shopping in stores will not wish to expend the mental effort in looking up conditions, noting down the most appropriate ones, and correlating this to particular products. However, such customers will spend the required time to merely select their characteristics if very little effort is required on their part, and if the apparatus operates in an interesting manner. A merchandising terminal which could aid customers in selecting the most appropriate products, in situations where the selection may depend on several factors or upon difficult-to-describe situations, would be of value in merchandising such products, which may include not only cosmetics, but other products as well."}
{"text": "1. Field\nExample embodiments disclosed herein relate to a display device to provide a three dimensional (3D) image and an internal image processing method, and more particularly, to an image processing apparatus and method that may decrease a visual fatigue when viewing a 3D image.\n2. Description of the Related Art\nRecently, interest in three dimensional (3D) image processing has increased. A 3D image may be configured by providing images corresponding to different views with respect to a plurality of views. For example, the 3D image may be a multi-view image corresponding to a plurality of views, or a stereoscopic image that may provide a left eye image and a right eye image corresponding to two views.\nCompared to a two dimensional (2D) display which displays only an existing plane image, a 3D display may provide a relatively higher sense of reality to a user. However, the 3D display may express different image information with respect to a left field of view and a right field of view and thus, may cause increased visual fatigue as compared to the 2D display. Accordingly, there is a desire for a solution to decrease visual fatigue.\nConventional methods for decreasing visual fatigue include a method of uniformly decreasing depth information, a method of increasing a contrast in the case of a near side region and decreasing the contrast in the case of a far side region, and the like."}
{"text": "Field of the Invention\nThe invention relates generally to the field of neutron generator-type well logging instruments. More specifically, the invention relates to structures for high voltage power supplies used with neutron generator-type well logging instruments to enable better placement of internal components of such instruments.\nBackground Art\nThe characteristics of geological formations are of significant interest in the exploration for, production and monitoring of subsurface water and mineral deposits, such as oil and gas. To that end, a variety of techniques have been developed to measure subsurface characteristics and evaluate the characteristics to determine certain petrophysical properties of interest. These techniques typically include the subsurface deployment of tools or instruments having energy sources to emit energy into the formations (usually from within a borehole traversing the formations). The emitted energy interacts with the surrounding formations to produce signals that are detected and measured by one or more sensors on the instrument. By processing the detected signal data, a profile or log of the subsurface properties is obtained.\nA variety of logging techniques have been developed to evaluate subsurface formations. A number of such techniques include emitting neutrons into the formation and evaluating the results of neutron interactions with formation nuclei. Various types of radiation sources have been used in subsurface well logging systems. For example, neutrons or gamma rays may be generated simply through the use of radioactive isotopes (which naturally decay over time), an x-ray source may be used or neutrons may be generated in an electronic device utilizing a nuclear fusion reaction to generate neutrons on demand. In such electronic neutron sources, high-energy neutrons may be generated through the controlled collision of energized particles by using a nuclear fusion reaction caused by accelerating ions toward a target at high voltage, so as to emit neutrons in bursts of fully controllable lengths and time sequences. Such devices may be referred to for convenience as neutron generators to distinguish them from chemical isotope sources. One neutron generator (referred to as a “pulsed” neutron generator) is described in U.S. Pat. No. 3,461,291. The neutron source described in the '291 patent uses an accelerator tube in which charged particles, such as deuterium ions, are accelerated across a high voltage potential and contact a target element such as tritium. The reaction between the deuterium ions with the tritium target produces substantially monoenergetic bursts of neutrons at an energy level of about 14 million electron volts (MeV). In most well logging applications the neutrons are not emitted continuously but in short bursts of well-defined lengths and sequence of repetition, however continuous generation of neutrons is also possible. When using such a neutron generator, the formation surrounding the instrument is typically subjected to repeated, discrete “bursts” of high energy neutrons. U.S. Pat. Nos. 4,501,964, 4,883,956, 4,926,044, 4,937,446, 4,972,082, 5,434,408, 5,105,080, 5,235,185, 5,539,225, and 5,608,215, for example, describe well logging instruments equipped with pulsed neutron generators.\nIn well logging using a neutron generator, the borehole and surrounding formation are irradiated with neutrons, and the various interactions of the neutrons with constituent nuclei are measured. Pulsed neutron well logging instruments typically include one or more sensors or detectors that record numbers of neutrons, particularly epithermal energy and thermal energy, as well as gamma rays which are emitted as a result of the interaction of the neutrons with the subsurface formations and the fluids in the borehole itself. The gamma rays may include inelastic gamma rays which are a consequence of high-energy collisions of the neutrons with atomic nuclei in the subsurface formations, as well as capture gamma rays emitted when low energy (thermal) neutrons are captured by susceptible atomic nuclei in the formations (for example, chlorine). Various relevant well logging techniques and tools are described, for example, in U.S. Pat. No. 4,390,783 to Grau, U.S. Pat. No. 4,507,554 to Hertzog et al., U.S. Pat. No. 5,021,653 to Roscoe et al., U.S. Pat. No. 5,081,351 to Roscoe et al., U.S. Pat. No. 5,097,123 to Grau et al., U.S. Pat. No. 5,237,594 to Carroll, and U.S. Pat. No. 5,521,378 to Roscoe et al.\nProperties of the formations which may be determined as a result of measuring neutron and gamma ray phenomena include, for example, formation density, fractional volume of void or pore space in the formation (porosity), carbon/oxygen (C/O) ratios, formation lithology, and neutron capture cross section (Sigma), among other measurements. Properties which may be determined by spectral analysis of the gamma rays include concentration of various chemical elements, for example. Properties of fluids in the wellbore may also be determined from various neutron and gamma ray measurements.\nNuclear measurements are also useable in nuclear spectroscopy techniques to obtain qualitative and quantitative information related to subsurface fluid movement. U.S. Pat. No. 5,219,518 describes an instrument equipped with a neutron source and sensors adapted to measure water flow through neutron oxygen activation. Alternative techniques for subsurface fluid measurements include the use of radioactive markers or tracers to identify flow path between formations or wells. U.S. Pat. Nos. 5,049,743, 5,182,051, 5,243,190, and 5,929,437 describe the use of elements that can be made radioactive by bombardment with neutrons so their location can be determined by nuclear logging. Logging tools equipped with gamma ray detectors are particularly suited to distinguish and determine the location of trace materials.\nThe nuclear phenomena detected with the foregoing instruments are representative of interactions not only with the formation nuclei, but also with the instrument and the borehole. In order to penetrate the formation, the high energy neutrons must pass through the fluid in the borehole (and casing in some applications) before entering the formation. The resulting non-formation contributions to the measured radiations significantly complicate the analysis of the formation characteristics. The problem is all the more complex since the sensitivity of the detector(s) to the radiations coming from the borehole, instrument and the formation, is a function of many parameters, such as, to name a few, lithology, porosity, borehole size, casing size/weight/eccentricity, cement quality, detector housings, or borehole fluid composition. In practice, several techniques have been devised to account for these contributions and to discriminate the undesired radiations from the desired radiations.\nFor certain types of neutron measurements, the neutron generator and its associated high voltage power supply may be disposed in an instrument housing along with neutron detectors, a neutron monitor detector, gamma ray detectors, or some combinations of the foregoing. Neutron generator well logging instruments known in the art include a high voltage power supply to operate the neutron generator disposed adjacent to the neutron generator. Certain of the foregoing detectors may need to be in nearly identical longitudinal positions as the neutron generator within the instrument in order to optimize the measurements made by the such detectors. Such optimal detector placement is impracticable using conventional instrument configurations where the neutron generator and its associated high voltage power supply are disposed proximate each other. FIG. 1 illustrates such conventional instrument configuration. The instrument 10 includes a spectroscopy detector 12 operable for gamma ray detection, neutron detectors at different longitudinal positions, e.g., near 14 and far/array 16, a pulsed neutron generator 18, and appropriate shielding 20.\nFIG. 2 illustrates another known instrument configuration. The instrument 30 is configured in what is called the “split-physics design”, in which smaller diameter detectors 32, 34 can be disposed longitudinally alongside the neutron generator 36. Larger detectors 38, such as the gamma-ray detector referred to above, can be longitudinally displaced from the neutron generator 36 so that the combination of detectors can be better optimized. Appropriate shielding 39 is also included in the instrument. The instrument shown in FIG. 2 is described, for example, in U.S. Pat. No. 7,148,471 assigned to the assignee of the present invention. Limitations inherent in the configuration of FIG. 2 includes that the size of the detectors 32, 34 near the neutron generator 36 must be relatively small, and/or the diameter of the instrument housing must be relatively large. The foregoing may limit the optimum placement of detectors that are typically larger in diameter, for example, spectral gamma ray detector, or may require the use of instrument housings of such size as to require the use of the instrument only in larger diameter wellbores.\nThe foregoing statements related to neutron well logging instruments may also apply to well logging instruments having x-ray generators therein.\nA need exists for improved radiation generator and detector configuration within a well logging instrument."}
{"text": "1. Field of the Invention\nThe present invention relates to an active-matrix-type display apparatus wherein pixel circuits each including light-emitting elements, such as organic EL (Electroluminescence) elements, are arranged in a matrix and a fabrication method for the display apparatus.\n2. Description of the Related Art\nIn an image display apparatus, such as a liquid crystal display unit, a great number of pixels are arrayed in a matrix, and in order to display an image, the light intensity is controlled for each of the pixels in response to the information of an image to be displayed.\nWhile this similarly applies to an organic EL display unit or a like display unit, an organic EL display unit is a self-luminous type display unit which has a light emitting element for each of the pixel circuits. The organic EL display unit is advantageous in that, in comparison with a liquid crystal display unit, the visibility of an image is high and a backlight need not be provided, and the speed of response is high.\nFurther, the organic EL display unit is much different from the liquid crystal display unit in that the gradation of color development is obtained by controlling the luminance of each light-emitting element through the current value applied to the light-emitting element. In short, the light-emitting element is of the current controlled type.\nWhile, in an organic EL display unit, a simple matrix system and an active matrix system can be applied as a driving system therefor, similarly to a liquid crystal display unit, the former system has a problem in that, while the structure is simple, it is difficult to implement a large display unit having high definition. Therefore, development of the active matrix system is carried out energetically, wherein current applied to a light-emitting element in the inside of each pixel circuit is controlled by an active element, usually by a TFT (Thin Film Transistor), provided in the pixel circuit.\nFIG. 1 shows a configuration of a general organic EL display apparatus.\nReferring to FIG. 1, the display apparatus 1 shown includes a pixel array section 2 wherein pixel circuits (PXLC) 2a are arranged in a m×n matrix, a horizontal selector (HSEL) 3, and a writing scanner (WSCN) 4. The display apparatus further includes signal lines or data lines SGL1 to SGLn selected by the horizontal selector 3 such that a data signal corresponding to luminance information is supplied thereto and scanning lines WSL1 to WSLm selectively driven by the writing scanner 4.\nIt is to be noted that the horizontal selector 3 and the writing scanner 4 may be formed from a MOSIC or the like on polycrystalline silicon or around the pixels MOSIC.\nFIG. 2 shows an example of a configuration of a pixel circuit 2a shown in FIG. 1. It is to be noted that the configuration is disclosed, for example, in U.S. Pat. No. 5,684,365 or Japanese Patent Laid-Open No. Hei 8-234683.\nReferring to FIG. 2, the pixel circuit shown has the simplest circuit configuration from among a great number of proposed circuit configurations and is a two-transistor driving type circuit.\nThe pixel circuit 2a includes a p-channel, thin-film field-effect transistor (hereinafter referred to as TFT) 11 and another TFT 12, a capacitor C11, and an organic EL light emitting element (OLED) 13 which serves as a light-emitting element. Further, in FIG. 2, reference characters SGL and WSL denote a signal line and a scanning line, respectively.\nSince an organic EL light-emitting element in most cases has a rectification characteristic, it is sometimes called an OLED (Organic Light Emitting Diode). While a symbol of a diode is used to indicate a light-emitting element in FIG. 2 and so forth, the rectification characteristic is not necessarily required for the OLED in the following description.\nThe TFT 11 is connected at the source thereof to a power supply potential Vcc, and the light-emitting element 13 is connected at the cathode thereof to the ground potential GND. Operation of the pixel circuit 2a shown in FIG. 2 is described below.\nStep ST1:\nThe scanning line WSL is placed into a selected state (here, into a low-level state), and a writing potential Vdata is applied to the signal line SGL. Consequently, the TFT 12 is rendered conducting to allow the capacitor C11 to be charged or discharged, and the gate potential of the TFT 11 is changed to the potential Vdata.\nStep ST2:\nThe scanning line WSL is placed into a non-selected state (here, into a high-level state). Consequently, the signal line SGL and the TFT 11 are electrically isolated from each other. However, the gate potential of the TFT 11 is retained in stability by the capacitor C11.\nStep ST3:\nThe current to be supplied to the TFT 11 and the light-emitting element 13 is changed to current which has a value corresponding to a gate-source voltage Vgs of the TFT 11, and the light-emitting element 13 continues to emit light with a luminance corresponding to the current value.\nThe operation for selecting the scanning line WSL to transmit the luminance information applied to a data line to the inside of a pixel, as at the step ST1, is hereinafter referred to as “writing.”\nAs described above, in the pixel circuit 2a shown in FIG. 2, once writing of the potential Vdata is performed, then the light-emitting element 13 continues to emit light with a fixed luminance until the next rewriting of a potential is performed.\nAs described above, in the pixel circuit 2a of FIG. 2, the gate application voltage of the TFT 11, which is a driving transistor, is changed to control the value of the current to be supplied to the EL light-emitting element 13.\nAt this time, the p-channel driving transistor is connected at the source thereof to the power supply potential Vcc, and the TFT 11 always operates in a saturation region. Therefore, the source of the driving transistor serves as a constant current source having a current value calculated in accordance with the following expression 1:Ids=½·μ(W/L)Cox(Vgs−|Vth|)2  (1)where μ, Cox, W, L, Vgs and Vth indicate the mobility of a carrier, the gate capacitance per unit area, the gate width, the gate length, the gate-source voltage of the TFT 11 and the threshold value of the TFT 11, respectively.\nIn the simple-matrix-type image display apparatus, each of light-emitting elements emits light only at a selected moment, but in the active-matrix-type image display apparatus, each light-emitting element continues to emit light also after the writing comes to an end, as described hereinabove. Therefore, the active-matrix-type image display apparatus is advantageous, particularly for a large display unit having a high definition, in that the peak luminance and the peak current of the light-emitting elements can be decreased in comparison with the simple-matrix-type image display apparatus.\nFIG. 3 illustrates the aged deterioration of a current-voltage (I-V) characteristic of an organic EL light-emitting element. Referring to FIG. 3, a curved line indicated by a solid line indicates a characteristic in an initial state, and a curved line indicated by a broken line indicates a characteristic after aged deterioration.\nGenerally, as times passes, the I-V characteristic of an organic EL light-emitting element deteriorates as seen in FIG. 3.\nHowever, in the two-transistor driving in FIG. 2, constant current is continuously supplied to the organic EL light emitting element in order to perform the constant current driving as described above. Therefore, even if the I-V characteristic of the organic EL light-emitting element deteriorates, the light emission luminance of the EL device does not deteriorate with time.\nIncidentally, while the pixel circuit 2a of FIG. 2 is formed from the p-channel TFTs, if the pixel circuit 2a can be formed from n-channel TFTs, then an existing amorphous silicon (a-Si) process can be used for TFT production. Consequently, a reduction of the cost of a TFT substrate can be anticipated.\nNow, a basic pixel circuit where each transistor is replaced with a n-channel TFT is described.\nFIG. 4 shows a pixel circuit wherein the p-channel TFTs in the circuit shown in FIG. 2 are replaced with n-channel transistors.\nThe pixel circuit 2b shown in FIG. 4 includes a n-channel TFT 21 and another n-channel TFT 22, a capacitor C21, and an organic EL light-emitting element (OLED) 23 serving as a light-emitting element. Further, in FIG. 4, reference characters SGL and WSL denote a data line and a scanning line, respectively.\nIn the pixel circuit 2b, the TFT 21 serving as a driving transistor is connected at the drain side thereof to the power supply potential Vcc and at the source thereof to the anode of the EL light-emitting element 23 in such a manner as to form a source follower circuit.\nFIG. 5 illustrates an operation point of the TFT 21 as the driving transistor and the EL light-emitting element 23 in an initial state. In FIG. 5, the axis of abscissa indicates the drain-source voltage Vds of the TFT 21, and the axis of ordinate indicates the drain-source current Ids of the TFT 21.\nAs seen in FIG. 5, the source voltage depends upon an operation point between the TFT 21 serving as the driving transistor and the EL light-emitting element 23, and has a value which differs depending upon the gate voltage.\nSince the TFT 21 is driven in a saturation region, the current Ids is supplied which has a current value calculated in accordance with the equation represented by the expression 1 regarding the voltage Vgs corresponding to the source voltage at the operation point."}
{"text": "1. Field\nThe present invention generally relates to error detection and correction mechanisms in computer memories. More specifically, the present invention relates to a technique that facilitates error detection and error correction after a failure of a memory component in a computer system.\n2. Related Art\nComputer systems routinely employ error-detecting and error-correcting codes to detect and/or correct various data errors which can be caused, for example, by noisy communication channels and unreliable storage media. Some error codes, such as SECDED Hamming codes, can be used to correct single-bit errors and detect double-bit errors. Other codes, which are based on Galois fields, can be used to correct a special class of multi-bit errors caused by a failure of an entire memory component. (For example, see U.S. Pat. No. 7,188,296, entitled “ECC for Component Failures Using Galois Fields,” by inventor Robert E. Cypher, filed 30 Oct. 2003, referred to as “the '296 patent.”) After a memory component fails, it is desirable to be able to detect and correct additional errors that arise during subsequent computer system operation. The technique described in the '296 patent can correct subsequent single-bit errors. However, this technique cannot be used to detect subsequent double-bit errors. It is also desirable to be able to reduce the number of additional “checkbits” which are used by this technique to provide such error correction and detection.\nHence, what is needed is a method and an apparatus for detecting and correcting errors that arise after a memory component has failed without the shortcomings of existing techniques."}
{"text": "This application is a 371 of PCT/EP98/04319, which was filed on Jul. 11, 1998.\nThe present invention relates to processes for the preparation of oxo products by the hydroformylation of substrates having Cxe2x95x90C double bonds using unmodified rhodium catalysts in a reaction mixture essentially consisting of the substrates, the catalyst and carbon dioxide in a supercritical state (scCO2). In particular, the invention relates to such processes for the preparation of products which contain substantial proportions of branched i-oxo products. Further, the invention relates to such processes for the hydroformylation of substrates which do not correspond to the general formula CnH2n. The invention further relates to such processes in which the separation of product and catalyst is effected using the special solvent properties of scCO2.\nThe transition metal catalyzed reaction of substrates containing Cxe2x95x90C double bonds with a mixture of hydrogen (H2) and carbon monoxide (CO) is referred to as hydroformylation or oxo reaction; it is a technically important method for the preparation of aldehydes and alcohols (oxo products). More than 6 million tons per year of oxo products are produced by catalytic hydroformylation worldwide. These products are employed as plasticizers and modifiers for PVC and other polymers, in detergent production, and as fine chemicals and as intermediates for the production of agrochemicals, food additives and pharmaceuticals (C. D. Frohning, C. W. Kohlpaintner, in Applied Homogeneous Catalysis with Organometallic Compounds (editors: B. Cornils, W. A. Herrmann), VCH, Weinheim, 1996, Vol. 1, Section 2.1.1).\nDepending on the reaction parameters, the catalysts and the substitution pattern at the Cxe2x95x90C double bond of the substrate, linear (normal, n-) or branched (iso, i-) oxo products are produced more or less selectively in the hydroformylation. Of the oxo products of simple short-chain olefins of general formula CnH2n (n=2, 3), it is mainly the n-products which are economically important. Of the long-chain olefins CnH2n (nxe2x89xa74), the branched i-oxo products are also economically important and are discussed, for example, as starting materials for the preparation of plasticizers in pure form or as mixtures with the n-oxo products. Above all, the i-oxo products are also particularly important for substrates with functionalized Cxe2x95x90C double bonds, i.e., substrates which do not correspond to the general formula CnH2n. In particular, reference may be made herein to the technical synthesis of vitamin A described in U.S. Pat. No. 3,840,589, DE 2 03 078 and U.S. Pat. No. 4,124,619, which involves the formation of an i-aldehyde by the hydroformylation of substrates comprising allyl ester moieties. The formation of branched products in the hydroformylation of vinyl aromatics (e.g., styrene) is discussed as a possible route to xcex1-aryl-carboxylic acids, which are employed, inter alia, as analgetics and antirheumatics (e.g., Ibuprofen(copyright), Naproxen(copyright), Suprofen(copyright)).\nThe catalysts for hydroformylation can be classified into so-called xe2x80x9cunmodifiedxe2x80x9d and xe2x80x9cmodifiedxe2x80x9d catalysts which are respectively preferred for particular processes or particular substrates according to the prior art (C. D. Frohning, C. W. Kohlpaintner, supra, p. 33ff). xe2x80x9cModifiedxe2x80x9d systems means catalysts in which the catalytically active metal component contains additional ligands, usually phosphorus compounds, in addition to H and CO for increasing the useful life and for controlling activity or selectivity. The term xe2x80x9cunmodified catalystsxe2x80x9d denotes all other metal compounds which are capable of forming catalytically active hydrido-carbonyl compounds under the reaction conditions. Technical importance has been achieved to date mainly by catalysts based on the metals cobalt (Co) and rhodium (Rh). The separation of the products and recovery of the catalysts is an important factor in the technical realization of hydroformylation reactions.\nSupercritical (xe2x80x9cfluidxe2x80x9d) carbon dioxide (scCO2), i.e., compressed carbon dioxide at temperatures and pressures beyond the critical point (Tc=31.0xc2x0 C., pc=73.75 atm, dc=0.467 gxc2x7mlxe2x88x921) is employed as a reaction medium for hydroformylation in U.S. Pat. No. 5,198,589 and in German Application DE-A-197 02 025.9 (Jan. 23, 1997). Carbon dioxide in a supercritical state is an interesting solvent for performing catalytic reactions because it is toxicologically and ecologically safe, in contrast to conventional organic solvents. Further, scCO2 has the property of being completely miscible with many gaseous reaction partners within wide limits, whereby limitation of the reaction rate by diffusion processes, which frequently occurs in gas/liquid phase reactions, is totally avoided. Further, due to the solvent properties of scCO2 which vary as a function of pressure and temperature, separation of main or side products from the reaction mixture is possible in favorable cases when the external parameters are appropriately selected. A survey of catalytic reactions in scCO2 is found in Science 1995, 269, 1065.\nIn DE-A-197 02 025.9, modified rhodium catalysts are employed for the hydroformylation in scCO2, specially developed phosphorus compounds ensuring high solubility of the catalysts in scCO2. Thus, the stated advantageous properties of scCO2 can be fully utilized with modified rhodium catalysts, but the preparation of the phosphorus ligands is an additional cost factor in possible technical applications. In U.S. Pat. No. 5,198,589, unmodified cobalt catalysts are employed for the hydroformylation of simple olefins CnH2n in scCO2. As compared to conventional solvents, similar reaction rates are obtained with significantly higher selectivities in favor of the linear n-oxo products. This increase in selectivity in favor of the n-oxo products is attributed to the use of scCO2 as the reaction medium.\nWe have now found that unmodified rhodium catalysts can be efficiently used for hydroformylation in scCO2 to obtain, surprisingly, not only significantly higher reaction rates, but also significantly higher selectivities in favor of the branched i-oxo products than are obtained in conventional solvents.\nRhodium catalyzed hydroformylations in scCO2 are conveniently performed by charging the catalyst or catalyst precursor and the substrate into a high-pressure reactor and then pressurizing with H2 and CO, either as a mixture or successively, at room temperature up to the desired partial pressure. Then, the amount of CO2 required to reach the desired density of the reaction medium is filled into the reactor. Then, heating to the desired reaction temperature is performed with stirring. After said temperature has been reached, the stirrer can be turned off due to the fast diffusion within the homogeneous supercritical phase. After the desired reaction time, the pressure is released from the reactor, during which the products can be isolated from the supercritical phase using appropriate known methods (K. Zosel, Angew. Chem. 1978, 90, 748; M. A. McHugh, V. J. Krukonis, Supercritical Fluid Extraction: Principle and Practice, Butterworths, Stoneham, 1994). Due to the fact that the solubilities of the metal species and the products are clearly different, the catalysts can be easily separated from the products and recycled.\nAs catalysts or catalyst precursors for the formation of unmodified Rh systems in scCO2, salts, complexes or cluster compounds of rhodium in any oxidation state can be employed. Compounds 1-8 are preferred examples of such catalysts or catalyst precursors, without intending to limit the choice of the catalyst to the structures shown. Particularly preferred catalysts are complex compounds of rhodium, 4-8, containing carbonyl ligands (CO) or ligands which can be readily replaced by CO under the reaction conditions. The amount of catalyst can be freely selected within a broad range depending on the reaction conditions and the reactivity of the substrates. Typical amounts of catalysts, based on the amount of substrate employed, range from 0.001 to 10 mole percent, preferably from 0.01 to 1 mole percent, more preferably from 0.05 to 0.5 mole percent. \nR1-R5 are residues independently selected from hydrogen, C1-C20 alkyl, C2-C20 alkenyl, C2-C20 alkynyl, aryl, C1-C20 perfluoroalkyl, C1-C20 carboxylate, C1-C20 alkoxy, C2-C20 alkenyloxy, C2-C20 alkynyloxy, aryloxy, C2-C20 alkoxycarbonyl, C1-C20 alkylthio, arylthio, C1-C20 alkylsilyl, arylsilyl, C1-C20 alkyloxysilyl or aryloxysilyl, each of which may optionally be substituted with C1-C12 alkyl, C1-C12 perfluoroalkyl, halo, C1-C5 alkoxy, C1-C12 carboxylate, C2-C12 alkoxycarbonyl or aryl. Residues R1-R5 may be connected with one another in cyclic compounds.\nL1-L4 are neutral ligands independently selected from the group consisting of CO, R1R2Cxe2x95x90CR3R4, R1Cxe2x89xa1CR2, NR1R2R3, wherein R1-R4 are as defined above. L1-L4 may be connected with one another in cyclic compounds.\nX, Y are anions or ligands with one negative charge, independently selected from the group consisting of Rxe2x88x92, Fxe2x88x92, Clxe2x88x92, Brxe2x88x92, Ixe2x88x92, ROxe2x88x92, RCO2xe2x88x92, p-C6H4SO3xe2x88x92, PF6xe2x88x92, BF4xe2x88x92, BR4xe2x88x92, wherein R is defined as mentioned for R1-R5.\nPossible substrates for the rhodium catalyzed hydroformylation in scCO2 are compounds containing at least one Cxe2x95x90C double bond which, on one hand, allow a hydroformylation reaction due to their substitution pattern, and on the other hand, are sufficiently soluble in scCO2 so that homogeneous solutions are obtained during the reaction. Preferred are substrates with a Cxe2x95x90C double bond as represented by general formula 9. Mixtures of such compounds may also be used as substrates. The amount of substrate employed is not critical as long as complete solubility in the supercritical medium is achieved. Based on the reactor volume, typical substrate quantities are 0.05-5 mol/l, preferably 0.1-2 mol/l. \nR6-R9 are residues independently selected from hydrogen, C1-C20 alkyl, C2-C20 alkenyl, C2-C20 alkynyl, aryl, C1-C20 perfluoroalkyl, C1-C20 carboxylate, C1-C20 alkoxy, C2-C20 alkenyloxy, C2-C20 alkynyloxy, aryloxy, C2-C20 alkoxycarbonyl, C1-C20 alkylthio, arylthio, C1-C20 alkylsilyl, arylsilyl, C1-C20 alkyloxysilyl, aryloxysilyl, C1-C20 alkylsilyloxy, C1-C20 alkyloxysilyloxy, C1-C20 alkylsulfonyl, C1-C20 alkylsulfinyl or halo, each of which may optionally be substituted with additional residues corresponding to the same definition as R6-R9. Residues R6-R9 may be connected with one another in cyclic compounds.\nThe partial pressure of H2 and CO can be widely varied. Typical pressures at room temperature are between 0.1 bar and 100 bar for each of H2 and CO, preferably between 1 bar and 30 bar. The influence of the H2 and CO partial pressure on the reaction rate and selectivity is largely identical with the trends found in conventional solvents (C. D. Frohning, C. W. Kohlpaintner, supra, p. 55ff). Much the same applies to the reaction temperature the lower limit of which is determined by the critical temperature of CO2 (Tc=31xc2x0 C.). Typical reaction temperatures for hydroformylation with unmodified rhodium catalysts in scCO2 are between 31xc2x0 C. and 150xc2x0 C., preferably between 35xc2x0 C. and 100xc2x0 C. When the H2/CO pressure and the reactor volume are given, the total pressure prevailing at the reaction temperature is determined by the amount of CO2. The minimum amount thereof is determined by the critical density of CO2 (dc=0,467 gxc2x7mlxe2x88x921), while the upper limit depends on the maximum admissible system test pressure of the reactors employed. For standard V2A high pressure reactors with inspection glasses, this results in a typical range for the amount of CO2 employed, based on the reactor volume, of from 0.46 gxc2x7mlxe2x88x921 to 0.90 gxc2x7mlxe2x88x921, preferably from 0.5 gxc2x7mlxe2x88x921 to 0.75 gxc2x7mlxe2x88x921.\nReactions according to the present invention may also be performed in the presence of one or more additives to provide, for example, an easier handling of the substrates or catalysts, or an improvement of the solvent properties of the reaction medium, or an increase of the reaction rate, or an improvement of the yield. Such additives may be independently selected, for example, from water, amines, perfluorinated compounds, organic solvents (e.g., dichloromethane, trichloromethane, tetrachloromethane, 1,2-dichloroethane, trichloroethene, benzene, toluene, xylene, cumene, hexane, cyclohexane, halobenzenes, tetrahydrofuran, tert-butyl methyl ether, diethyl ether, dimethoxyethane, dimethylformamide, ethyl acetoacetate, acetone, dimethyl carbonate, alcohols).\nIn order to ensure that the reaction will start only after the supercritical state is reached, it may be convenient, depending on the reactivity of the substrate, to supply the catalyst and the substrate in a spatially separated way. Thus, for example, either the substrate or the catalyst may be supplied to the reactor in a separate storage vessel open at the top. One component may also be supplied in a second pressure vessel which is connected with the reactor through a valve. It is also possible to supply the catalyst in glas ampoules which can be cracked under the reaction conditions. Another alternative is the continuous or discontinuous metering of a component by means of a pump system.\nExample 1 describes the typical procedure adopted in the hydroformylation in scCO2 with unmodified rhodium catalysts and illustrates the increased reaction rate and the unexpected change of selectivity when scCO2 is used as compared to a conventional solvent. The Examples summarized in Table 1 describe the hydroformylation of some prototypical compounds in scCO2 under typical conditions, but are not intended to limit the scope, the application range or the advantages of the present invention in any way."}
{"text": "The tracking of high value assets is an ongoing business problem for many companies. In particular, the tremendous volume of goods flowing through a transportation system on any given day creates a significant logistical challenge. As a result, tracking and visibility systems play an integral part in most transportation systems by allowing both customers and transportation personnel to track the flow of goods. However, despite the existence of many prior art systems, items flowing through a transportation system are, on occasion, still delivered to a location later than expected, misdelivered to an incorrect address, and even lost or stolen somewhere along the way. These and other such mishaps can be particularly troubling when the items being shipped are considered “high-value” items, based on either their relative importance to the shipper or on their raw economic value. Examples of high-value items can include shipments of diamonds, expensive jewelry, computer chips, cars or sensitive documents, to name just a few.\nFor the most part, item tracking services known in the art are passive in the sense that they are generally not designed to actively address the types of problems discussed above. In other words, most solutions in the past have been inhibited by their lack of active scanning and reporting of location for each particular item moving through the transportation network. This is due in part to the fact that many items are bundled or aggregated together with other items for efficiency of shipment. As a result, many items are not physically (i.e., actually) scanned at every transportation point, which they pass through. Instead, a “logical” scan of the item is performed, meaning that if the bundled or aggregated load arrives at (or departs from) a particular location, the bundled load is scanned and the tracking data for each particular item assumed to be within that bundle is updated accordingly. This occurs despite the fact that each particular item believed to be in the bundle may not actually be present.\nThus, for example, if an item is lost, stolen, behind schedule, or misdelivered to an incorrect address, such problems may go undetected until a shipper or consignee notices the problem and calls to report it. By then, however, a substantial amount of time may have been passed, during which transportation and security personnel could have been looking into the problem. As such, the likelihood of being able to diagnose and remedy the cause of any potential problems may have been reduced.\nFurthermore, while it is sometimes possible to ship an item via special shipping channels that incorporate more stringent security measures, doing so may not always be desirable. Aside from the fact that shipping an item in such a way can be substantially more expensive, the added security measures can, in some cases, create a new set of problems by drawing unnecessary attention to the item. In other words, by shipping the item via special shipping channels, the shipper may in effect be identifying the item to would-be criminals as a high value item. Thus, in some cases it may actually be safer to at least provide the appearance that the item is being shipped via “standard” means so that the item blends in and, as such, does not stand out as an item of particular value in the various shipping yards, hub locations and other transport points through which the item is likely to pass.\nTherefore, an unsatisfied need exists in the industry for improved tracking and visibility systems and methods that overcome the deficiencies in the prior art, some of which are discussed above."}
{"text": "Electronic systems generally include at least one printed circuit board (PCB) containing one or more integrated circuit (IC) chips or ICs. ICs typically include input/output (I/O) pins which may be coupled to various interconnects of the PCB. Testing performance of electronic systems which include PCBs and ICs typically requires testing at multiple levels including at the chip level, at the board level, and at the system level. Testing at the board level includes testing interconnects of the PCB. Testing at the system level requires analysis of interconnections between and among the ICs, the PCBs, and other devices both on and off the PCB.\nTo enhance testability at the board level as well as at the system level, a common design practice at the chip level is to incorporate boundary scan test logic into an IC in accordance with IEEE Standard 1149.1. 1149.1 specifies the function of boundary scan logic known as JTAG, which is named for the Joint Test Action Group, for control of boundary scan testing. Two basic elements of an IC are a core logic and the I/O pins. In accordance with 1149.1, boundary scan cells (BSCs) are inserted between the core logic and the I/O pins of the IC. BSCs are typically inserted for all I/O pins of the plurality of ICs on the PCB and may be used to test the integrity of the interconnections between the plurality of ICs.\nEach IC may be controlled by boundary scan logic, in accordance with 1149.1, to operate either in a system mode or in a JTAG test mode. In the system mode, system data signals relating to core functions of the IC are passed through the I/O pins to and from devices external to the IC. In the JTAG test mode, test data are provided by the boundary scan chain for the purpose of testing interconnections between the IC and devices external to the IC. The boundary scan test access port (TAP) controller also provides test control signals which include mode signals, shift signals, clock signals, and update signals, among others, each of which is well known. The mandated public instructions include a bypass instruction, a sample instruction, a preload instruction, and a extest instruction. The extest instruction controls BSCs to perform a boundary scan test among the various ICs.\nThe IC further includes a test data input (TDI) demultiplexer, a test data output (TDO) multiplexer, a test access port (TAP) controller, and a plurality of test registers such as a bypass register, an instruction register, and an identification register. The TDI demultiplexer includes an input coupled to receive a test data signal from the boundary scan logic which is typically driven externally to the IC. The TDI demultiplexer includes a first output coupled to a TDI input of a first BSC of the plurality of BSCs in the IC. Each of the BSCs includes a TDI input and a TDO output. Each of BSCs is connected serially from a TDO output to a TDI input to propagate test data signals from one BSC to the next BSC in the chain. The TDI demultiplexer further includes a second output coupled to an input of the core logic, a third output coupled to an input of the bypass register; a fourth output coupled to an input of the instruction register; and a fifth output coupled to an input of the identification register.\nThe TDO multiplexer includes an output which is coupled to provide a test data signal to another IC or to the boundary scan logic. The TDO multiplexer further includes: a first input coupled to a TDO output of a last BSC of the plurality of BSCs in the IC, a second input coupled to an output of the bypass register; a third input coupled to an output of the instruction register, and a fourth input coupled to an output of the identification register. The identification register includes inputs coupled to outputs of the TAP controller. The TAP controller includes inputs coupled to receive a TMS signal, a TCK signal, and a TRST signal from the boundary scan logic.\nIn general, there are three possible I/O structures for an IC including a two-state I/O structure, a three-state I/O structure, and a bi-directional I/O structure. Each of the three I/O structures provides coupling between the core logic and at least one I/O pin. Any or all of the I/O structures may be used in an IC depending on the particular circumstances. The two-state I/O structure includes a two-state output buffer having an input and an output. The input of the two-state output buffer is coupled to a system data output of the core logic. The output of the two-state output buffer is coupled to an I/O pin. The three-state I/O structure includes a three-state output buffer having an input, an output, and a control input. The input of the three-state output buffer is coupled to a system data output of the core logic. The output of the three-state output buffer is coupled to an I/O pin. The control input of the three-state output buffer is coupled to a three-state system control signal output line of the core logic. The bi-directional I/O structure includes a bi-directional buffer. The bi-directional buffer includes an output buffer element having an input, an output, and a control input and an input buffer element having an input and an output. The control input of the output buffer element is coupled to a bi-directional control signal output line of the core logic. The input of the output buffer element is coupled to a system data output of the core logic. The output of the input buffer element is coupled to a system data received input of the core logic. The output of the output buffer element and the input of the input buffer element are coupled together with an I/O pin.\nAccording to conventional methods and apparatus for boundary scan testing, the BSCs are inserted into the I/O structures between the buffers and the core logic. For a two-state output structure, a BSC is inserted between the core logic and the input of the two-state output buffer. For a three-state output structure, a BSC is inserted between the system data output of the core logic and the input of the three-state output buffer. Also, a BSC is inserted between the three-state control signal output line of the core logic and the control input of the three-state output buffer. For a bi-directional output structure, a BSC is inserted between the system control signal output line of the core logic and the bi-directional output buffer. Also, a bi-directional BSC is inserted between the core logic and the bidirectional output buffer.\nIEEE Standard 1149.1 was first adopted in 1990. It has been widely used and has proved to be very successful. 1149.1 has been amended twice to improve it. However, 1149.1 does not address all situations and design practices. One such practice is the inclusion of capacitive coupling in the interconnections between ICs. A capacitor is added either to the connection between the ICs or to one, the other, or both of the I/O pins of the ICs or the PCBs with connectors. The capacitor is designed to reduce noise and block unwanted common mode voltage differences in the interconnection. For discussion, this will be referred to alternatively as either being AC coupled or DC de-coupled.\nTurning first to FIG. 1, a block diagram of ten possible combinations of DC and AC coupled interconnections between two devices is shown. The choice of which of the combinations shown that are actually used depends on the circumstances. Because of the capacitor, the value of a signal at the receiving end of the interconnection is no longer the same as the value at the driving end. This assumes that the steady state condition has been reached where all transient values of driven data have subsided. The square wave input signal is transformed at the receiving end into a series of decaying signal spikes corresponding to the transitions of the square wave. The rate of the spike decay depends on the value of the capacitor and the inherent resistance, that is, the RC constant. The result is that conventional 1149.1 testing becomes impractical on AC coupled interconnections. One will note that there are seven possible AC coupled combinations where 1149.1 will not work as compared to only three DC coupled combinations where 1149.1 will work. As the quest for higher signal speeds continues in the future, the use of AC coupling will increase. This becomes especially true with the development of optical communication signals. The consequence will be less and less reliance on conventional 1149.1 testing.\nBeyond the direct connected AC coupled testing are the areas of AC coupled cluster testing and AC coupled functional block testing. Under certain circumstances, one might want to test a cluster or functional block of components. As these components might not be directly connected to one another, further testing difficulties might arise. For example, the components might have multiple AC coupled connections in series or they might have non-electrical connections. If so, the test signal propagation may require continuous streams of alternating patterns in order to maintain minimum cutoff frequency requirements of the signal path or to maintain prolonged test signal application to compensate for long signal delays. Delay is one of the difficulties that might be encountered and is a measure of the time that a signal takes to propagate from the transmitter to the receiver. Another way to express the delay of the signal is in terms of the latency of the communication path. Generally the delay is constant but initially unknown. Another potential difficulty is pin skew which is a measure of the variation of two or more signals that are designed to be simultaneous in theory but are not quite in practice. Generally the skew is constant but initially unknown. Yet another potential difficulty is jitter which is a measure of the variation of a signal over time, that is, a signal that is supposed to have uniform timing in theory is not quite in practice. Jitter can be especially problematic when it arises in one or more of the clock signals as a consistent clock signal can be important for accurate timing. Generally the jitter is within bounds that are not initially known. It is not strictly necessary that each of these three difficulties be precisely measured in order to compensate for them, but they should be accounted for to produce accurate testing results. A flexible approach is proposed."}
{"text": "In the field of integrated circuits, the use of plastic molded packages for providing mechanical and environmental protection for the integrated circuit chips has become commonplace. Such plastic molded packages are primarily beneficial due to their low cost of manufacture relative to ceramic encapsulated packages, particularly in avoiding the high cost of the piece parts of ceramic encapsulated packages.\nComplementary metal-oxide-semiconductor (CMOS) technology for the design and fabrication of integrated circuits has also become more popular in recent years, due to the very low level of standby power dissipation in circuits fabricated according to this technology. Low power CMOS integrated circuits have enabled the widespread use of battery power in modern electronic systems, whether as backup or as primary power, as the operating life of the batteries in powering these low power circuits can be quite long.\nImplementation of battery-powered integrated circuits into electronic systems, even as backup power, requires an additional degree of complexity in the system design and manufacture. As such, it has been desirable for some time to integrally package an electrical cell, or battery, within the integrated circuit package. Copending application Ser. No. 07/995,665, filed Dec. 21, 1992, entitled \"Surface Mountable Integrated Circuit Package with Integrated Battery Mount\" assigned to SGS-Thomson Microelectronics, Inc. and incorporated herein by this reference, describes an integrated circuit package of the surface-mountable type within which a battery is mounted to leads that extend from the side of an internal package body, with a housing attached to the package over the battery. Gaps between the housing and the internal package body provide thermal insulation of the battery from the package body, so that the circuit may be subjected to solder reflow mounting to a circuit board in the conventional sense for surface mountable packages, while insulating the battery from the high soldering temperatures.\nBy way of further background, it is desirable and customary in the art to test each integrated circuit in a fully functional manner, so that the integrated circuit customer can receive lots of the packaged circuit that are substantially free of defects. In many cases, however, integrated circuits having a battery backup mode cannot have their backup functionality tested with the actual battery with which they will be installed, but instead can only be tested in a circumstantial manner with an external battery source. Furthermore, while prior hybrid package techniques that include a battery with the integrated circuit are known, such prior techniques require high cost ceramic piece parts (such as headers, lids and the like), and thus are not suitable for use with high quantity commodity circuits such as random access memories, which are to be sold in highly price-competitive markets.\nIt is therefore an object of the present invention to provide a molded integrated circuit package within which both an integrated circuit and a battery power source are molded.\nIt is a further object of the present invention to provide such a package which enables complete functional testing of the integrated circuit, including the functionality of its battery backup capability.\nIt is a further object of the present invention to provide such a package which includes an inner package for the integrated circuit chip itself, so that assembly of the circuit and battery package may be done with a reduced risk that the integrated circuit chip will be damaged.\nOther objects and advantages of the present invention will be apparent to those of ordinary skill in the art having reference to the following specification together with the drawings."}
{"text": "The invention relates to laboratory equipment for chemical laboratories and more specifically to vacuum filter funnels of the shape typically referred to as a buchner funnel.\nA buchner funnel typically has an upper cylindrical portion which holds the liquid and solids to be filtered which upper portion terminates in a lower perforated plate. A filter medium such as filter paper is placed on the upper surface of the perforated plate and typically a vacuum is drawn on the lower side of the perforated plate to increase the flow rate of filtrate through the filter medium.\nIf the filter medium does not conform closely to the perforated plate, it is possible that solids will pass around the outside of the filter medium and through the relatively large perforations in the bottom plate. Such bypassing can be highly detrimental to many chemical laboratory procedures. Although this problem does not occur with frequency with common filter paper, it can be a substantial problem with rigid filter materials such as fritted glass."}
{"text": "1. Field of the Invention\nThe present invention relates to a biocide composition, more particularly to a biocide composition having synergistic effects by comprising 3-isothiazolone and polyhexamethyleneguanidine phosphate and a method for sterilizing microorganisms, fungi, etc. using the same.\n2. Description of the Related Art\nMicroorganisms, bacteria, mold, algae, etc. grow in industrial water such as cooling water, wastewater, emulsifiers used in the textile industry, etc. and have a deteriorating effective on the operations of industrial processes. Such microorganisms propagate using organisms contained in the industrial water as a nutrient source and secrete polysaccharides. Varieties of organic and inorganic substances combined with these secreted polysaccharides and form viscous lumps or masses also called slime. Organic materials such as cellulose, hemicellulose, and the fibrin of white water in paper industry in particular provide to have sufficient nutrient sources for such microorganisms. The slime formed at areas of low fluid flow in a paper processing process results in both direct and indirect losses, such as manufacturing time loss, equipment efficiency deterioration, etc. due to the deterioration of pulp quality, etc. Furthermore, the growth of microorganisms at such places like a cooling water facility, where much water is contained or recirculated causes a fouling phenomena. This deteriorates heat transfer efficiency in an industrial cooling tower, as well as corrodes metal or erodes wooden parts.\nBacteria such as taloblastic prokaryotes can propagate by decomposing various types of organisms, some of which are the sources of the microbiologically induced corrosion, by secreting and discharging polysaccharides forming a biofilm. Fungi such as eukaryote can also propagate by decomposing various types of organisms like bacteria, and some types of fungi secreting cellulase are the sources for discoloration and decomposition of wooden parts by way of degrading the fibroid materials of wooden parts of a cooling tower, etc.\nAlgae, such as eukaryote, can propagate by photosynthesizing under an environment with light, air, and a small amount of organisms, and carbohydrate formed from algae is used as a nutrient source for other microorganisms such as bacteria and fungi, thereby accelerating the fouling phenomena. Algal fouling due to algae propagation is intensified at places that are exposed to sunlight, particularly in cooling water facilities, swimming pools, etc. This phenomena results in clogging which blocks water pipes, as well as the deterioration of heat transfer efficiency, oxidization of metal surfaces by the generation of oxygen, and the promotion of corrosion at a holes on metallic surfaces through a partial galvanic reaction when the organism dies.\nVarious biocides are being developed in order to kill such microorganisms, fungi, algae, etc. or to prevent their surface adhesion to metal, etc. These biocides are generally divided into oxidant biocides and non-oxidant biocides. The oxidant biocides that are mainly used are halogen compounds such as chlorine, bromine, etc. that are popular due to their economic advantage of strong oxidation capability and low price. However, they can cause erosion of the wooden parts of cooling towers and metallic decomposition, and their sterilizing efficiency tends to deteriorate as they can be easily being discharged into atmosphere. Furthermore, their practical effectiveness is poor due their peculiar way of first reacting with secreted polysaccharides before they react with microorganisms that are the actual sources of biofilm formation.\nNon-oxidant biocides which overcome these disadvantages include 3-isothiazolone, quaternary ammonium salt, formaldehyde emission compound, glutaraldehyde, etc. and are mainly used separately. Although 3-isothiazolone, which is disclosed in U.S. Pat. Nos. 3,761,488, 4,105,431, 4,279,762, etc., has a high sterilizing effect and wide antibiotic spectrum, it has a disadvantage in that its immediate instantaneous sterilizing effects are low. Furthermore, Korean Patent Application No. 89-20381 discloses an antiseptic composition useful in preventing circulation water putrefaction wherein a biocides 5-chloro-2-methyl-4-isothiazolone-3-on and 2-methyl-4-isothiazolone-3-on are mixed in a ratio of about 3:1 and wherein this antiseptic composition further comprises didecyldimethylammoniumchloride. U.S. Pat. No. 4,379,137 discloses a method for improving sterilizing capability by mixing polymer quaternary ammonium salt and 3-isothiazolone. However, since these mixtures emit corrosive materials, i.e., halogen compounds such as fluorine, chlorine, etc., it is difficult for them to be applied where metals susceptible to corrosion are used, such as carbon steel, cast iron, stainless steel, copper, etc.\nAdditionally, Korean Patent Application No. 97-80170 discloses a process that can be applied even when metals susceptible to corrosion are used since the disclosed process does not emit a halogen compound. Another water soluble biocide composition having synergistic effects comprising 3-isothiazolone, which has high sterilizing effect on quartemary ammonium phosphate and microorganisms, is disclosed in Korean Patent Application No. 97-46517 as a biocide having superior properties of immediate sterilizing effects on microorganisms, durability, anticorrosiveness, etc. However, there is a problem in applying this invention to various industrial fields due to the issues such as bubbling or foaming, etc. when quaternary ammonium is added to provide immediate effectiveness to the 3-isothiazolone.\nFurthermore, while polyhexamethyleneguanidine phosphate has immediate effectiveness, and is used in the effective and wide control of microorganisms in various industrial fields, including water treatment, as well as exhibiting low foaming properties, it has a disadvantage of not having a wide antibiotic spectrum by itself."}
{"text": "This invention relates generally to light-emitting semi-conductor devices, or light-emitting diodes (LEDs) according to more common parlance, and particularly to those which lend themselves to use, by being arrayed, in small displays, now known as microdisplays, and in the printheads of LED printers, among other applications. The invention also specifically pertains to a method of fabricating such LED arrays.\nMicrodisplays and printheads incorporating an array of LEDs are now winning ever-increasing commercial acceptance. Generally, an LED array may be envisaged as a group of LEDs arranged in rows and columns on a substrate of electrically insulating or highly resistive material. The individual LEDs are electrically interconnected by conductors in a required pattern for individual excitation.\nAmong the materials that have been used for LED array substrates are sapphire, silicon carbide, and silicon. Both sapphire and silicon carbide permit layers of semiconducting compounds to be grown favorably thereon. Offsetting this benefit is their expensiveness, being considerably more costly than silicon. An additional drawback of sapphire and silicon carbide is that they are both more permeable to light, blue in particular, than is silicon. The light radiated by the LEDs toward the substrate may penetrate the same and hit its bottom, thereby to be diffusely reflected back through the neighboring LEDs (i.e., dots or pixels) and hence to leak out.\nThe silicon substrate has no such shortcomings. Besides being cheap, it is free from light leakage. Silicon is less in permeability to light than is sapphire or silicon carbide, so much so that the light from the LEDs is mostly absorbed by the substrate, with little or no leakage through the neighboring LEDs.\nJapanese Unexamined Patent Publication No. 2004-195946 is hereby cited as dealing with an LED array on a high-resistance silicon substrate. The LEDs are each constituted of a first semiconductor layer of, say, n-type gallium arsenide (GaAs) and, thereover, a second semiconductor layer of, say, p-type GaAs for emission of light in response to voltage application to these constituent layers. The LEDs glow individually by having their first semiconductor layers electrically interconnected in rows or columns by first wires bonded to their surfaces, and their second semiconductor layers electrically interconnected in columns or rows by second wires bonded to their surfaces. The first semiconductor layers partly protrude from under the second semiconductor layers to provide ledges to which the first wires are bonded.\nAn objection to this known LED array is that the first semiconductor layers with their protruding ledges must themselves provide parts of the current paths to the LEDs. In order for the first semiconductor layers to perform this function to the full, they must be much thicker than in the case where they are coupled directly to the first wires or equivalent parts, without any such protruding parts. The thicker first semiconductor layers make the complete device correspondingly thicker. Additionally, the thicker first semiconductor layers add to the costs of the materials needed for their fabrication and to the lengths of time required for their epitaxial growth. Thus the prior art LED array was unnecessarily bulky and costly.\nA further inconvenience with the prior art manifested itself when the first semiconductor layers of the LEDs were connected to one drive terminal. Such connection required the second wires to be branched at several points, necessitating highly complex LED wirings."}
{"text": "Social networking has become a popular medium for online collaboration. Maturing Web 2.0 technologies have influenced competing social networking sites (for example, Orkut, Facebook, etc.) to look at ways to provide innovative features. Also, telecommunications (telecom) technologies have experienced changes such as, for example, content and capabilities owned by the operator are now available for developers to combine. Additionally, social networking now includes the use of mobile devices to enable communications with on-line social networking buddies.\nHowever, traditionally, social networks have worked on relatively static and/or slow-changing information, that is, defined once by a user and changed infrequently. Groups that are defined in a social network in existing approaches are similarly static and pre-defined in terms of membership. Further, ways of communicating with group members are also statically defined (for example, leave a scrap, write on a wall, etc.)."}
{"text": "Schizophrenia is a chronic debilitating mental illness affecting about one percent of the population. The disease manifests in delusional behavior, dysfunctional thinking, agitated body movement, social withdrawal, and depression. Schizophrenia patients suffer a profoundly reduced quality of life, and are ten times more likely to commit suicide that the general population.\nDopamine (particularly D2 and D3) antagonists are well recognized as improving symptoms of schizophrenia, and have been used clinically as such for decades. In the past twenty years it has become recognized that treatment of schizophrenia, as with many mental illnesses, benefits from engaging multiple receptors including serotonergic and adrenergic. Despite, literally, dozens of approved drugs to treat schizophrenia the disease remains poorly treated in many patients. Side effects of current medications include: dyskinesia, akathisia, weight gain, mood disturbances, sexual dysfunction, sedation, orthostatic hypotension, hypersalivation, and (in some cases) arganulocytosis.\nAmisulpride (4-amino-N-(((1-ethyl-2-pyrrolidinyl)methyl)-5-(ethylsulfonyl))-2-methoxybenzamide) is an antipsychotic patented in 1981. Amisulpride binds selectively to the human dopaminergic D2 (Ki 2.8 nM) and D3 (Ki 3.2 nM) receptor subtypes without any affinity for D1, D4 and D5 receptor subtypes. Unlike classical and atypical neuroleptics, amisulpride displays low affinity for serotonin, alpha-adrenergic, histamine receptor subtypes, muscarinic receptors and sigma sites though it has also been demonstrated to bind 5-HT2B and HT7a receptors with low double digit nM Ki. This ability of amisulpride to bind 5-HT receptors is thought to result in amisulpride's ability to treat symptoms of depression (sometimes noted in schizophrenia patients). Interestingly, compared to other antipsychotics, amisulpride is not noted to have any activity at the 5-HT2a receptor.\nDespite the unique activities of amisulpride, amisulpride has low ability to cross blood brain barrier (BBB) to interact with the receptors in the brain. In a 2014 study, passive diffusion of amisulpride across a PAMPA membrane (Pe) was the lowest of 30 psychiatric drugs tested. Thus, dosing of amisulpride is high, typically 400 to 800 mg/d (though up to 1,200 mg/day is not uncommon). Such a high dose may cause adverse effects to the treated subjects."}
{"text": "Iminosugars are potent inhibitors of glycosidases. Azasugars of the isofagomine family are inhibitors of configuration-retaining β-glycosidases due to the formation of a strong electrostatic interaction between a protonated endocyclic nitrogen at the anomeric center of the imino sugar and the catalytic nucleophile of the enzyme. The inhibitors mimic the transition state in the hydrolysis of the glycosidic bond. Isofagomine, (3R,4R,5R)-3,4-dihydroxy-5-hydroxymethylpiperidine, also known as IFG, is one glycosidase inhibitor which was synthesized in anticipation that it would be effective as a liver glycogen phosphorylase inhibitor for the treatment of diabetes (see U.S. Pat. Nos. 5,844,102 to Sierks et al., and 5,863,903 to Lundgren et al., both of which are herein incorporated by reference).\nIFG Tartate Salt, its production and its use to treat Gaucher Disease has also been described in U.S. patent application Ser. No. 11/752,658, which is hereby incorporated by reference."}
{"text": "1. Field of the Invention\nThis invention relates to a light shielding mechanism for a quadrilateral opening and more particularly to a light shield mechanism for such an opening within a single lens reflex camera and the like.\n2. Description of the Prior Art\nIn single lens reflex cameras, provided with a mechanism for automatically determining the aperture size or the extent of the exposure time or similar parameters in accordance with photographing conditions such as the brightness of the field to be photographed, the camera employs an internal light receiving photometric system in which a photosensitive portion for measuring the light within the field to be photographed arranged within the region of the light path emanating from the field to be photographed and transmited through an image lens in order to determine the aforementioned parameters for film exposure. However, it is not possible to disregard the influence of counter-incident light which falls on the light path from the eyepiece portion of the view finder optical system. Particularly where the eye piece portion is exposed when the photograph is taken by means of a self-timer or remote control mechanism, a relatively large error is produced by the presence of the counter-incident light acting on the light being measured relative to the field being photographed.\nFor this reason, it is often required that a mechanism be provided for shielding the view finder eyepiece in single lens reflex cameras and the like, and it is further necessary to add such a mechanism for shielding a light projection optical system when a slide change is being effected within a slide projector to thereby enhance projection effects, and similarly it is often required that such a shield mechanism be provided within optical illumination systems in telescopic photograhic apparatus or the like.\nIt is, therefore, an object of the present invention to provide a light shield mechanism for a camera view finder eyepiece and the like which reduces the area of the shield blades incorporated within said mechanism where the optical path across which the light shield is employed comprises a quadrilateral opening such as a rectangle or the like.\nIt is a further object of the present invention to provide an improved shield mechanism which minimizes the space required for the shield blades within the periphery about the quadrilateral light path opening and to provide a light sield mechanism which is both simple in construction and which readily attains the desired light shielding function."}
{"text": "1. Field of the Invention\nThe present invention relates to an echo canceler, which cancels echoes generated during hands free speaking on a cellular telephone.\n2. Description of the Related Art\nEchoes generated during the hands free speaking mean that a speech reception voice generated from a speaker of communication equipment is re-input to a microphone and is transmitted as echoes to a far end talker through a transmission circuit.\nTherefore, it is required to prevent from transmitting the speech reception voice through the transmission circuit back to the far end talker at receiving equipment. To purse the requirement, it is necessary to develop an echo canceling technique.\nThe Japanese laid-open patent application No. 7-297901 describes as one of conventional techniques, voice communication equipment with an echo canceler to cancel echoes generated during the hands free talking.\nIt is possible to cancel the echoes by estimating and subtracting echoes included in a speech transmission signal from a speech reception voice. In the echo canceling technique, described in the Japanese patent publication No. 7-297901, a variable attenuator is activated to fix the transmission signal level at the hands free speaking but not during the both far end and near end concurrently speak so that noise level is sent and unnatural feeling is given to the far end talker.\nFurther, the other technique relating to an echo canceler to cancel the echoes generated during the hands free speaking is described in the Japanese laid-open patent application No. 6-216811. It is an object of the technique is to solve a problem such that distortion of the speech voice signal occurs when a speech reception voice signal, of which level is large, is input, thus the echo signal can not be canceled sufficiently at a speech transmission signal path.\nTo achieve the object, in the description of the Japanese laid-open patent application No. 6-216811, the speech reception voice signal level is limited to a predetermined amplitude level or less before the voice signal is supplied to an echo canceler and a speaker. However, there are some cases where the limitation of signal level gives an influence not preferable for the voice signal output from the speaker.\nAccordingly, it is an object of the present invention to provide an echo canceler to improve preferable echo canceling, which is not improved sufficiently by the conventional techniques.\nIt is another object of the present invention to provide an echo canceler whereby a speech reception voice power can be adaptably controlled and the distortion of voice output from a speaker can be reduced.\nIt is further object of the present invention to provide an echo canceler to improve a stable echo canceling operation by attenuating a voice input to the echo canceler, varying the attenuation amount, and removing the distortion of the voice data.\nIt is furthermore object of the present invention to provide an echo canceler for general use, which does not depend on the delay amount and gain of an echo pass by automatically measuring echo pass delay, echo pass gain and clipping state in an analog circuit.\nTo achieve the above-described objects according to the present invention, an echo canceler comprising, a power control section controlling a power for a speech reception voice, a first filtering section removing a particular frequency component from an output of the power control section, an echo canceling section removing echo components of a speech reception voice, which are added to a speech transmission voice, a second filtering section removing the particular frequency components of the transmission voices which include echo, and a state judgement section judging existences of the speech reception and transmission voices from an output of the second filtering section and controlling a power attenuation amount of the power control section based on the judgement result.\nIn one preferred mode of the present invention, the state judgement section calculates the result of judging whether or not the speech reception and transmission voices exist and a power of the speech reception voice of N samples, based on N samples of the speech reception voice and the output from the second filtering section, and there is comprised of a delay buffer maintaining the calculated result of the judgement section and the power of the speech reception voice for T blocks, each block having N samples.\nIn the above-described structure, the state judgement section controls to make the power control of the power control section operable, when both the speech reception and transmission voices exist.\nIn the other mode of the present invention, the state judgement section makes the attenuation amount obtained by the power control of the power control section at the time when either of the speech reception or transmission voice does not exist smaller than that at the time when both of the speech reception and transmission voices exist.\nAdditionally, the echo canceling section estimates a echo replica according to an attenuated signal obtained by attenuating an input signal to the echo canceling section, amplifies the estimated echo replica according to the attenuation of the input signal, and subtracts the amplified echo replica amplified from the input signal."}
{"text": "1. Field of the Invention\nThe present invention relates generally to air conditioning systems and, more specifically, to an accumulator insulator bracket for an air conditioning system in a vehicle.\n2. Description of the Related Art\nIt is known to provide an accumulator for an air conditioning system in a motor vehicle. The accumulator provides for the accumulation of liquid and liquid/vapor separation of a working fluid within the air conditioning system. Typically, an accumulator includes a housing having a first section and a second section, that are joined together to define an interior chamber. The housing also has an inlet opening through which a refrigerant may be introduced into the interior chamber, and an outlet opening through which refrigerant may exit the interior chamber. The accumulator also includes a separate inlet tube, outlet tube, oil return and desiccant bag disposed in the interior chamber of the housing. An example of an accumulator is disclosed in U.S. Pat. No. 5,729,998 to Grohs et al., issued Mar. 24, 1998, the disclosure of which is incorporated by reference.\nPreviously, the accumulator was attached to a portion of the air conditioning system, such as the evaporator case, using a bracket. An example of a bracket is a metal strap that wraps around the accumulator and is secured to a structure using a threaded fastener. While this type of bracket works well, a disadvantage is that it is difficult to consistently position the bracket during the vehicle assembly process.\nIt is known that the performance of the air conditioning system is affected by temperature variations of the liquid or liquid/vapor within the accumulator. Therefore, the accumulator is typically surrounded by an insulating device to insulate the accumulator from environmental temperature conditions, and also prevent heat transfer to the accumulator. An example of an insulating device is a foam sleeve having open ends. While the foam sleeve prevents a certain amount of heating, a disadvantage is a lower insulating efficiency since the accumulator is not completely encased. Another example of an insulting device is disclosed in U.S. Pat. No. 6,041,618 to Patel et al. entitled xe2x80x9cInsulated Pressure Vessel Holderxe2x80x9d. This insulating device also has a sleeve shape, as well as a corrugated interior wall and an integral mounting bracket. In addition to the disadvantage of having an open-ended sleeve shape, a further disadvantage of this insulating device is that the corrugated shape of the sleeve encourages thermal contact between the sleeve and the accumulator, resulting in a decreased insulating value. Thus, there is a need in the art for an accumulator insulator bracket that efficiently and cost effectively insulates and supports the accumulator for an air conditioning system.\nAccordingly, the present invention is an accumulator bracket for an accumulator in an air conditioning system. The accumulator bracket includes a housing having a cylindrical shape with a first closed end, a second closed end and a wall extending between the first closed end and the second closed end. The housing is longitudinally divided into a first housing member and a second housing member. The accumulator insulating bracket also includes an air flow directing rib extending radially from and axially along an interior portion of the wall of the housing, such that the rib defines an insulating air pocket between the wall and the accumulator and supports the accumulator.\nOne advantage of the present invention is that an accumulator insulator bracket is provided that integrates attachment and insulation of the accumulator into one structure. Another advantage of the present invention is that the accumulator insulator bracket has a clamshell shape to fully enclose the accumulator to prevent heat transfer. Still another advantage of the present invention is that an inside wall of the accumulator insulator bracket includes integrally molded-ribs optimally arranged to create air pockets around the outside of the accumulator that insulate the accumulator. A further advantage of the present invention is that the accumulator insulator bracket can be integrally molded with a housing for another component of the air conditioning system, such as the evaporator.\nOther features and advantages of the present invention will be readily appreciated, as the same becomes better understood, after reading the subsequent description taken in conjunction with the accompanying drawings."}
{"text": "Recently, an image-recording material has been dominated particularly by a material for forming a color image and specifically, a recording material of inkjet system, a recording material of heat-sensitive transfer system, a recording material of electrophotographic system, a silver halide photosensitive material of transfer system, a printing ink, a recording pen and the like are popularly used. Further, in an image device such as CCD of photographing equipment or in a display such as LCD and PDP, a color filter is used for recording and reproducing a color image. In these color image recording materials and color filters, colorants (dyes or pigments) using three primary colors for a so-called additive or subtractive color mixing method are used for displaying or recording a full color image, but a colorant having absorption characteristics capable of realizing a preferred color reproduction region and having fastness enough to endure various use conditions or environmental conditions has not been found until now, and accordingly, improvements are keenly demanded.\nThe dyes and pigments used in each of the aforementioned uses need to commonly have the following properties. That is, the dyes and pigments need to have preferred absorption characteristics for color reproducibility, and exhibit good fastness under environmental conditions used, for example, light resistance, heat resistance, resistance to oxidative gases such as ozone and the like. In addition, in the case where the colorant is a pigment, the colorant also needs to have properties such as being substantially insoluble in water or an organic solvent, showing good chemical resistance, and not impairing the preferred absorption characteristics in the molecular dispersion state even when used as a particle. The requisite characteristics may be controlled by varying the degree of the intermolecular interaction, but both characteristics are in a trade-off relationship and therefore, it is difficult to satisfy both at the same time.\nFurther, in using a pigment, other than the properties described above, the pigment also needs to have a particle diameter and a particle shape necessary for bringing out a desired transparency, to show good fastness under environment conditions used, for example, light resistance, heat resistance, resistance to an oxidative gas such as ozone, water resistance, and chemical resistance to an organic solvent, a sulfurous acid gas or the like, and to have properties capable of being dispersed even into microparticles in a medium used and keeping the dispersed state stable.\nThat is, compared to the dye which is required to have performances as a colorant molecule, performances required for the pigment need to satisfy not only performances as a colorant molecule, but also the aforementioned requisite performances as a solid (fine particle dispersion) as an aggregate of colorant molecules in various fields. As a result, a compound group which may be used as a pigment is extremely limited as compared to the dye, and thus even when a high-performance dye is derived into a pigment, a pigment capable of satisfying the requisite performances as a fine particle dispersion is very few in number and may not be easily developed. This is also confirmed from the fact that the number of pigments registered in the Color Index is less than 1/10 of the number of dyes.\nAmong the pigments, particularly an azo pigment has high lightness and excellent light resistance and heat resistance, and thus is widely used as a pigment for a printing ink, an inkjet ink, an electrophotographic material, and a color filter. Further, as the use is expanded, the pigment also needs to have much better stability over time regardless of a medium used rather than stability at a level usually used in a printing ink, a gravure ink, and a coloring agent. For example, Patent Document 1 discloses an azo pigment with a specific structure, which has excellent color characteristics such as hue and excellent light resistance and includes a pyrazole ring.\nSimultaneously, in the use of a color filter or an inkjet ink, it is required to further enhance clearness or transparency. In order to enhance clearness or transparency, it is effective to finely disperse a pigment, and therefore, a method of efficiently preparing an azo pigment fine particle by which a fine dispersion may be formed is also required. Patent Document 2 discloses a method of preparing an azo pigment with a specific structure, which includes a pyrazole ring, by which the azo pigment may be prepared at a high efficiency and a low cost.\nMeanwhile, most of the representative organic pigments are polymorphic, and it is known that these pigments take two or more crystal forms in spite of having the same chemical composition.\nIn the organic pigments, for example, there are also an organic pigment which may obtain fine and size-regulated particles by selecting appropriate reaction conditions during synthesis like the azo pigment, an organic pigment obtained by preparing a pigment from very finely and aggregated particles produced during synthesis by subjecting the particles to particle growth and size-regulating the particles in subsequent steps like a copper phthalocyanine green pigment, and an organic pigment obtained by preparing a pigment from coarse and uneven particles produced during synthesis by finely pulverizing and size-regulating the particles in subsequent steps like a copper phthalocyanine blue pigment. For example, a diketopyrrolopyrrole pigment is generally synthesized by reacting succinic diester with aromatic nitrile in an organic solvent (see, for example, Patent Document 3). Then, the crude diketopyrrolopyrrole pigment is heat-treated in water or an organic solvent, followed by pulverization such as wet grinding into a form suitable for use (see, for example, Patent Document 4). In the C.I. Pigment Red 254, an α-crystal form and a β-crystal form are known (see, for example, Patent Document 5). In addition, in the C.I. Pigment Yellow 181 which is an azo pigment, various types of crystal forms are known (see, for example, Patent Document 6)."}
{"text": "1. Field of the Invention\nThis invention relates generally to minimally invasive medical implants and procedures for their use and insertion device for implanting the same, having particular application for treating urinary incontinence.\n2. Description of the Prior Art\nStress urinary incontinence (SUI) is a female medical condition commonly associated with weakening of the pelvic muscles and/or connective tissues that support the urethra in its proper position. As a result of this condition, involuntary urine leakage occurs from simple physical activity, such as running or jumping, and even coughing or sneezing, as the urethra is not properly supported and does not remain fully closed during such activity.\nA widely accepted medical procedure to correct SUI is the insertion of a tension-free or trans-vaginal tape that is surgically implanted in the pelvic tissue and that extends under and provides support for the urethra when pressure is exerted thereon.\nU.S. Pat. No. 5,899,909, the disclosure of which is incorporated herein by reference, describes in detail a typical procedure for treating SUI using a trans-vaginal tape. The tape is implanted by passing an elongated curved needle that is attached to one end of the tape through an incision in the vaginal wall, to one lateral side of the urethra, through the pelvic tissue behind the pubic bone, and exiting out through an incision made in the abdominal wall. The procedure is repeated for the other end of the mesh tape, this time on the other lateral side of the urethra, with the needle exiting through a second incision made in the abdominal wall of the patient. After the mesh tape is adjusted for proper support of the urethra, its free ends extending outside of the abdominal wall are trimmed. Over time, fibroblasts grow into the mesh tape to anchor the tape in the surrounding tissue. Thus, the tape is left as an implant in the body to form an artificial ligament supporting the urethra in order to restore urinary continence. In another known method for implanting a trans-vaginal tape, the tape is inserted in a somewhat similar manner, but is brought out through the obturator hole and exits the body through a small incision in the upper leg.\nThe use of trans-vaginal tape for treating SUI has a number of advantages. It does not need to be attached through bone anchors, sutures or any other element to secure the tape in place, and there is minimal scarring. The procedure takes about 30 to 50 minutes, and may be performed on an outpatient basis under local, regional or general anesthesia. One of the few disadvantages of known procedures for implanting sub-urethral tapes is that the use of needles to pass the tape through the body poses a risk for vessel, bladder and bowel perforation. Also they require two separate, minimal incisions made through the abdominal wall (for a retropubic approach) or the upper leg (for an obtuator approach) through which exit the curved needles and attached tape is passed. This, of course, increases the risk of post-operative pain and/or infection to at least a small degree.\nAccordingly, what is needed is an improved sub-urethral tape, and device and method for implanting the same."}
{"text": "The present invention relates in general to an apparatus and method for injecting faults on the pins of an integrated circuit employed to perform a logical function.\nIn the testing of the error detection section of a device employing integrated circuits, it is desirable to be able to create error conditions on the inputs and outputs of the integrated circuits employed in the device. The ability to create such error conditions is also useful in the development and testing of diagnostic programs used in data processing systems to diagnose faults in the data processor and attached input/output devices.\nTesting apparatuses for dynamically injecting errors in a data processing system have previously been proposed, for example, in East, D. G., et al., Error Injector for Testing a Data Processing Unit, in IBM Tech. Disc. Bull., 17 (6): p. 1691, Nov. 1974, and also in Maiden, D. W., et al., Error Injector Probe, in IBM Tech. Disc. Bull., 18 (1): p. 7, June 1975. In such apparatuses, it will be realized that an error may only be injected at one point in a device once the testing apparatus is connected. Further, the occurrence and location of an error cannot be controlled programatically.\nComputer controlled circuit testers for applying predetermined signals to individual pins to which the circuit to be tested is connected have also been proposed, for example, in U.S. Pat. No. 3,872,441 by George William Cailow, filed Nov. 29, 1972. In such an apparatus, it will be realized that errors cannot be applied to circuits that are performing a logical function in the device of which they are an element.\nTesting units for insertion in a circuit with a circuit board, including an extender board containing circuits for transmitting and receiving data, have been proposed, for example, in U.S. Pat. No. 3,904,861 by John E. McNamara, filed Mar. 13, 1974. In such an apparatus, the circuit board is removed and the extender board is placed in circuit with the circuit board. The circuit board may then be tested with a duplication of system conditions, the circuits in the extender board receiving and transmitting data. In such an apparatus, it will be realized that the functions performed by the circuits on the extender board are fixed by the hardware and therefore signals cannot be selectively altered or passed through without modification. Further, the apparatus is designed to operate only on signals entering and leaving the circuit board through the circuit board's connection with the connector block the circuit board normally plugs into. As a consequence, no modification to signals passing between circuits mounted on the circuit board can be performed.\nTherefore, it is an object of the present invention to provide an error injection apparatus which can simulate error conditions in an integrated circuit which is an in-line component of an electronic circuit.\nAnother object of the present invention is to provide an error injection apparatus capable of simulating errors on multiple input or output pins of an integrated circuit.\nIt is a further object of the present invention to provide an error injection apparatus capable of varying the input or output pins of an in-line integrated circuit which have simulated errors applied to them without the necessity of rewiring the error injection apparatus.\nIt is an additional object of the present invention to provide a programmatically controlled apparatus and method for simulating errors on input or output pins of an integrated circuit while the integrated circuit is performing its normal logical functions as an in-line component of an electronic circuit."}
{"text": "1. Field of the Invention\nThe invention relates to an apparatus for storing and transporting a spare tire. More particularly, the invention relates to a spare tire carrier which is removable from the vehicle and doubles as a hand dolly for transporting the tire to and from the vehicle.\n2. Description of the Prior Art\nCurrently spare tires for trucks and multiple wheeled vehicles are stored in a tire carrier rack placed under the vehicle. To retrieve the spare tire from the rack the tire must be pulled from the rack and slid or rolled along the ground. Storing the spare tire in the rack, similarly, requires sliding or rolling the spare tire under the truck and lifting it into the rack. Considering the weight and size of tires used on multiple wheeled vehicles, the operation described above is physically difficult and generally requires that the operator obtain assistance.\nSeveral patents have issued which unsuccessfully address this problem. U.S. Pat. No. 2,489,561 issued to Kenneth C. Clark, on Nov. 26, 1946, describes an apparatus in which a tire is mounted on a swingable arm which is in turn mounted on the frame of the vehicle. The swingable arm allows the tire to be stored under the vehicle and allows the arm to pivot and place the tire in a vertical position near the axis for tire mounting operations. U.S. Pat. No. 3,554,397 issued to John Verl Cluff, on Nov. 29, 1968, describes an apparatus which is useful for flat bed vehicles. The spare tire is connected directly by a cable to two pullies and a winch. The winch is turned and the tire is pulled along the ground and lifted into a position under the vehicle frame."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for the carbonization of oil shale and/or other carbonizable materials, including a process line serving for carbonization, with a closed cycle and without condensate recovery, wherein the materials flowing through the process line are heated to the carbonization temperature through heat exchange.\nComing into consideration as further carbonizable materials are fuels which are high in inerts such as, for example, bituminous rock, oil-containing fuller's earth, high-ash coal, oil sands, refuse and the like.\n2. Discussion of the Prior Art\nNumerous processes have already been proposed for the carbonization of oil shale and similar fuel-containing materials having a high proportion of inert materials, for some of which pilot plants have also been constructed. However, it is difficult to completely utilize the fuel content of the raw material and with a satisfactory degree of efficiency since, in essence, subsequent to the carbonization there remains a predetermined proportion of fuel, particularly carbon, in the carbonized material, whose heat of combustion can frequently not be employed, and if so then only with extensive and complex apparatus.\nA survey of such processes is elucidated in E/MJ (September 1977) pages 148 through 154. One of the previously mentioned processes is described therein under the designation \"Lurgi L.R.\" In the known process the decarbonized raw material which is conducted through the cycle is employed as a heat carrier. After the carbonization this raw material is pneumatically conveyed into a storage bin, whereby there is burned the so-called fixed carbon so as to lead to an increase in the sensible heat beyond the carbonization temperature. The conveying air, which is concurrently utilized as combustion air for the combustion of the fixed carbon, leaves this process line into the atmosphere at combustion temperature and is lost. Also the heat which is contained in the ash, in effect, in the fully burned-out raw material, is lost.\nIn another process which is referred to in that publication, ceramic spheres are heated through the combustion of process gas, and which serve as solid heat carriers for the heat which is required during the carbonization process. For the heating of the ceramic spheres there is utilized a special sphere heater whose exhaust gases serve for the preheating of the raw material which is to be carbonized such as oil shale. In this process, the fixed carbon which still remains in the raw material subsequent to the carbonization is not utilized. The carbonized residue leaves the process at the carbonization temperature (approximately 500.degree. C.) and must be cooled by being sprayed with water since it will otherwise continue to burn and represent a significant burden on the environment (for instance, a stench).\nOther processes which are mentioned therein utilize carbonization gas which is conveyed in a closed cycle as a heat carrier. This gas is heated either directly through combustion or through indirect heating. Employed as the cycle gas is the incondensible component of the carbonization gases. Thus, the cycle gases are cooled to condensate precipitation, so as to be thereafter again heated for heat transfer for the carbonization process. This results in a significant heat loss with a correspondingly lowered degree of efficiency. A still satisfactory degree of heat is possessed by a process wherein the heat which is transferred to the cycle gas originates from the combustion of the fixed carbon. Even in this process does the ash leave the process at a combustion temperature of about 800.degree. C.\nMoreover, there have also been proposed retort processes with indirect heat transfer from the gas to the solid material, however, these have not been found any applicability since the indirect heat transfer of gas/solids material in the required large operating units, due to the large heat transfer surfaces then required, lead to presently uncontrollable constructions.\nAdditionally, there have also been previously proposed so-called on-site processes in which the raw material is carbonized in its original deposit whereby the fixed carbon should also presently be burned. These on-site processes have also failed to find application in view of the difficulty soluable problems of a controlled process conductance. For such processes for relatively high-yield oil shale, reports have been made of a yield of merely about 50 to 60%.\nIn the decarbonization of other raw materials, processes have already been proposed which are not suitable for the carbonization of carbonizable material such as oil shale and similar raw materials.\nFor example, from German Pat. No. 11 60 823 there has become known a process for continuous degassing, such as the carbonization and/or coking of fine-granular, non-caking water-containing fuels through the intermediary of hot gas streams. This process is particularly suited for the production of coke as the primary product, whereas carbonization gas occurs as a side-product which is thinned-out through the combustion gases of the heat generator. For the carbonization of oil shale and other carbonizable raw materials with high contents of inert materials and with an also necessarily low fuel content; however, this process is not suitable since the process residue still contains the fixed carbon and no combustion is provided. Since the residue exists from the process at carbonization temperature, for a large quantity of residue the heat requirement of the process is extremely high whereby so as to also require extremely large quantities of combustion gas in order to be able to maintain the necessary process temperature. The thereby obtained thinning of the carbonization gases reduces the heat value of the noncondensable gases so extensively for introduced low-grade materials, that it practically precludes any utilization of the gases. For fuel-poor, ash-rich materials, such as oil shale or the like, whose commerical utilization is the object of the invention, this process is also not suitable. The same also pertains to a similar known process as well as an arrangement for the continuous thermal treatment, such as degasification and/or coking, of fine-granular water-containing fuels; having reference to German Laid-Open Patent Application 16 71 320. Also in this instance it is impossible to recover an unthinned high-quality carbonization gas. The carbonization gas extensively contains the exhaust gases of an external heat generator.\nIn another known process for the degasification of powdered fuels (as in German Pat. No. 977,218) there is employed a distillation chamber in which the material is loosened by means of the introduction of CO.sub.2 or steam to a density of 8 to 80 g/l. The introduced gases thin out the product gas. The required heat quantity is conveyed into the distillation chamber through the intermediary of a solid heat carrier from a further similar treatment chamber, and in which there is combusted the remaining carbon and, as occasioned, auxiliary fuel. The excess combustion residue and the exhaust gases from this combustion stage leave the process at the temperature of the combustion stage. The thusly generated high heat losses can only by avoided by extensive additional measures for facilitating utilization of the exhaust gas heat.\nFinally, also known is a process and an arrangement for the degasification of fine-granular up to dust-like fuel (refer to German Pat. No. 11 97 432), in which there are employed as essential components at least two air heaters acting as heat exchangers as primarily employed in the foundry industry. This process operates in an alternating manner whereby one of the air heaters is currently heated for the combustion of the incondensible gas and the stored heat of currently the other air heater is transferred to a mixture consisting of carrier gas and the fuel which is to be degasified. This process is also not suited to the carbonization of oil shale and similar carbonizable fuels which are high in inerts, since the fixed carbon material cannot be utilized and, in essence, because the air heaters are heated with incondensible gas. Moreover, the carbonized residue is removed at the carbonization temperature so that the heat losses are also appreciable."}
{"text": "1. Field of the Invention\nThe present disclosure relates to methods of treating light-colored wood substrates to produce a wood substrate with a low gloss, open pore appearance that provides enhanced protection from ultraviolet (“UV”) radiation and a “real” wood look. The present disclosure also relates to the wood components produced by the methods disclosed herein.\n2. Description of the Related Art\nReal wood trim pieces with very glossy finishes in some automotive vehicles are difficult to distinguish from plastic woodgrain films. A simple polyurethane (herein “PU”) coating can have the desired appearance but typically cannot meet rigorous performance targets for, for example, long term resistance to fading upon extended exposure to sunlight.\nHowever, real wood remains a popular choice for the interior of upscale and luxury automobiles. In order to meet rigorous performance targets multiple coats of high gloss finishes are routinely used with real wood trim pieces in automotive vehicles in an attempt to meet the light resistance standard and other performance criteria, but the resulting high gloss finish makes the real wood appear to the ordinary consumer to be less desirable woodgrain plastic films.\nA method of achieving the necessary light resistance properties while maintaining a “real” wood appearance, that is, not a high gloss finish, is of interest."}
{"text": "1. Field of the Invention\nThe present disclosure relates to the field of fabrication of microstructures, and, more particularly, to a tool for chemically mechanically polishing substrates, bearing for instance a plurality of dies for forming integrated circuits, wherein the tool is equipped with a slurry delivering unit for delivering slurry onto the surface of a polishing pad of the tool.\n2. Description of the Related Art\nIn microstructures, such as integrated circuits, a large number of elements, e.g., transistors, capacitors and resistors, are fabricated on a single substrate by depositing semi-conductive, conductive and insulating material layers and patterning those layers by photolithography and etch techniques. Frequently, the problem arises that the patterning of a subsequent material layer is adversely affected by a pronounced topography of the previously formed material layers. Moreover, the fabrication of microstructures often requires the removal of excess material of a previously deposited material layer. For example, individual circuit elements may be electrically connected by means of metal lines that are embedded in a dielectric, thereby forming what is usually referred to as a metallization layer. In modern integrated circuits, a plurality of such metallization layers is typically provided, wherein the layers are stacked on top of each other to maintain the required functionality. The repeated patterning of material layers, however, creates an increasingly non-planar surface topography, which may cause deterioration of subsequent patterning processes, especially for microstructures including features with minimum dimensions in the submicron range, as is the case for sophisticated integrated circuits.\nIt has thus turned out to be necessary to planarize the surface of the substrate between the formation of specific subsequent layers. A planar surface of the substrate is desirable for various reasons, one of them being the limited optical depth of the focus in photolithography which is used to pattern the material layers of microstructures.\nChemical mechanical polishing (CMP) is an appropriate and widely used process to remove excess material and to achieve global planarization of a substrate. In the CMP process, a wafer is mounted on an appropriately formed carrier, a so-called polishing head, and the carrier is moved relative to a polishing pad while the wafer is in contact with the polishing pad. A slurry is supplied to the polishing pad during the CMP process and contains a chemical compound reacting with the material or materials of the layer to be planarized by, for example, converting into a reaction product that may be less stable and easier removed, while the reaction product, such as a metal oxide, is then mechanically removed with abrasives contained in the slurry and/or the polishing pad. To obtain a required removal rate while at the same time achieving a high degree of planarity of the layer, parameters and conditions of the CMP process must be appropriately chosen, thereby considering factors such as construction of the polishing pad, type of slurry, pressure applied to the wafer while moving relative to the polishing pad and the relative velocity between the wafer and the polishing pad. The removal rate further significantly depends on the temperature of the slurry, affected by the amount of friction created by the relative motion of the polishing pad and the wafer, the degree of saturation of the slurry with ablated particles and, in particular, the state of the polishing surface of the polishing pad.\nMost polishing pads are formed of a cellular microstructure polymer material having numerous voids which are filled with slurry during operation. A densification of the slurry within the voids occurs due to the absorbed particles that have been removed from the substrate surface and accumulated in the slurry. As a consequence, the removal rate steadily decreases, thereby disadvantageously affecting the reliability of the planarizing process and thus reducing yield and reliability of the completed semiconductor devices.\nTo partly overcome this problem, a so-called pad conditioner is typically used that “reconditions” the polishing surface of the polishing pad. The pad conditioner includes a conditioning surface that may be comprised of a variety of materials, e.g., diamond that is embedded in a resistant material. In such cases, the exhausted surface of the pad is ablated and/or reworked by the relatively hard material of the pad conditioner once the removal rate is assessed to be too low. In other cases, as in sophisticated CMP apparatus, the pad conditioner is continuously in contact with the polishing pad while the substrate is polished.\nIn modern integrated circuits, process requirements concerning uniformity of the CMP process are very strict so that the state of the polishing pad has to be maintained as constant as possible over the entire area of a single substrate as well as for the processing of as many substrates as possible. Consequently, the pad conditioners are usually provided with a drive assembly and a control unit that allow the pad conditioner, that is, at least a carrier including the conditioning surface, to be moved with respect to the polishing head and the polishing pad to rework the polishing pad substantially uniformly while avoiding interference with the movement of the polishing head. Therefore, one or more electric motors are typically provided in the conditioner drive assembly to rotate and/or sweep the conditioning surface suitably.\nOne problem with conventional CMP systems resides in the fact that a wafer removal profile, as well as a wafer removal rate, depends on many factors, e.g., the type of slurry, the slurry thickness, the temperature of the slurry, the pressure applied to the wafer while moving relative to the polishing pad, the relative velocity between the wafer and the polishing pad and the curvature of the wafer. Controlling a conventional CMP system therefore requires complex controlling of multiple parameters. Moreover, the deterioration of one or more of the consumables of a CMP renders it difficult to maintain process stability and to reliably predict an optimum time point for consumable replacement. Generally, replacing the consumables at an early stage significantly contributes to the cost of ownership and a reduced tool availability, whereas a replacement in a very advanced stage of the consumables of a CMP system may jeopardize process stability.\nThe present disclosure is directed to various systems and methods that may avoid, or at least reduce, the effects of one or more of the problems identified above."}
{"text": "As is well known, and widely accepted, partial task simulators and training aids can be very effective for teaching individuals how to perform a wide variety of different tasks. More specifically, they can be extremely helpful for teaching an individual how to perform certain medical procedures. In this context, and of particular importance for the present invention, are those medical procedures that are required for response to a life-threatening, emergency situation. The import here is two-fold. Firstly, the partial task simulator should effectively augment the educational background that is necessary to assess an emergency situation. Secondly, it should serve as a tool with which a person can learn how to respond to an emergency situation by properly performing essential life-saving tasks. The efficacy of any partial task simulator or training aid, however, is dependent on the realism it provides and its ability to simulate or mimic an environment where the task is to be actually performed.\nWith the above in mind, a catastrophic event presents a situation wherein the proper training of emergency medical personnel can be invaluable. Regardless whether the event is the result of an accident, a natural disaster or some form of combat, the consequence of a first response to the event may make the difference between life and death. In such instances, the ability of medical personnel to rapidly and reliably attend to wounds and injuries is of crucial importance. Practice on partial task simulators such as medical mannequins, while valuable as teaching aids, are limited by the mannequin's immobility, weight, expense and minimal interaction with the medical personnel.\nIn light of the above, it is an object of the present invention to provide a device for realistically and dynamically simulating the wounds and injuries on a person (e.g., role player, actor) that can be received during a traumatic event. Another object of the present invention is to provide a device that effectively functions as a training aid to teach a person how to treat the wounds and injuries that can be received by a person during a traumatic event. Another object of the present invention is to provide a device that effectively functions as a training aid that allows verbal and gesticular interaction between a live human wearing the device and a first responder who is treating the person wearing the device. Still another object of the present invention is to provide a training aid for teaching how to treat wounds and injuries that is easy to use, is simple to manufacture and is comparatively cost effective."}
{"text": "Accuracy and repeatability are required for any reliable test method. This is particularly important for testing of medical devices in which an error in testing can result in failing to administer the desired dosage of a therapeutic agent. One such device is a wearable medication pump.\nA wearable medication pump, such as an insulin pump, is a small device about the size of a small cell phone that is worn externally. The wearable medication pump delivers precise doses of a therapeutic agent to a patient by injection, and can be adjusted manually and/or through feedback sensors to closely match the patient's needs. The therapeutic agent is typically held within a reservoir installed in a reservoir compartment in the wearable medication pump. A pump motor slide advances a piston in the reservoir to deliver the therapeutic agent to the patient.\nBecause the dose of therapeutic agent delivered is a function of pump motor slide position, testing of the wearable medication pump must assure that the pump motor slide does not move unnecessarily and only moves the distance requested to deliver the desired dose. One method of measuring pump motor slide position is to insert a ruler into the reservoir compartment until the ruler contacts the pump motor slide, then note the distance from a point on the wearable medication pump casing to the point where the ruler contacts the pump motor slide. Another method of measuring pump motor slide position is to position a caliper against the pump casing, advance the caliper depth probe into the reservoir compartment until the depth probe contacts the pump motor slide, then read the distance on the caliper.\nUnfortunately, such distance measurement methods are neither accurate nor repeatable. The person performing the test can place the ruler or caliper depth probe into the reservoir compartment at different angles. When reading the distance on the ruler, the person performing the tests can use different landmarks points on the wearable medication pump casing. The test results can vary in a single test and can vary even more over a number of tests when different persons are likely to be performing the test. This, in turn, prevents achievement of highest product quality, product reliability, and patient safety.\nIt would be desirable to have a medication pump test devices and methods of use that would overcome the above disadvantages."}
{"text": "An image forming apparatus for forming an image on a recording sheet includes a developing device for supplying a developing agent to a surface of a photosensitive element. The developing device is generally configured to have a developing agent storage chamber and a developing chamber that can communicate with each other via a supplying port. A developing agent supplying roller and a developing roller are accommodated in the developing chamber.\nThe image forming apparatus, in some cases, has to dispose the developing device therein such that the developing agent storage chamber is located above the developing chamber, as disclosed, for example, in JP-A-2001-166556. The image forming apparatus disclosed in JP-A-2001-166556 adopts a horizontal tandem-type image forming system in which plural photosensitive drums are aligned in a horizontal direction and in parallel to one another to form a color image on a recording sheet by developing agents. The horizontal tandem-type image forming system suffers from a problem in that the length becomes long in a direction in which plural photosensitive drums are aligned.\nTo solve this problem, the developing devices which correspond, respectively, to the photosensitive drums have to reduce their dimensions in the direction in which the photosensitive drums are aligned. To this end, the developing agent storage chamber, which is normally located adjacent to the developing chamber, is located above the developing chamber as disclosed in JP-A-2001-166556.\nThe developing device disclosed in JP-A-2001-166556 can use gravity to downwardly move the developing agent from the developing agent storage chamber to the developing chamber, thereby supplying the developing agent to the feeding roller and the developing roller.\nHowever, since the developing agent is only moved downwardly from the developing agent storage chamber to the developing chamber in a one-way direction, repetition of development results in increase in a ratio of deteriorated developing agent in the developing chamber. In particular, by fresh developing agent being additionally supplied and forced into the developing chamber downwardly from the developing agent storage chamber, the deteriorated developing agent is accumulated at longitudinal end portions of the developing roller in the developing chamber. The accumulated developing agent is likely to lower the quality of an image considerably at both end portions of the image, in particular, when a color image is formed.\nNot only the horizontal tandem-type color laser printer but also a monochrome laser printer suffers from the same problem when a developing agent storage chamber is located above a developing chamber.\nSince the problem of deterioration of a developing agent is caused by subjecting the developing agent to sliding friction, the problem commonly can occur over developing devices in which a developing agent is charged through friction charging. In the event that the developing agent is deteriorated, phenomena such as a reduction in fluidity due to the embedment of externally added agents which are externally added to the developing agent and a reduction in charging quantity due to a reduction in charging function are generated. This phenomena may cause rough and scattered printing and in worst case, background fogging, on an image formed on a recording sheet.\nWhen a fresh and non-deteriorated developing agent is added to the deteriorated developing agent whose charging function is reduced, a charging failure is caused in the deteriorated developing agent depending upon conditions, and this may cause a situation in which the background fogging is worsened."}
{"text": "This invention relates to gas producing and handling systems and, in particular, to such systems which produce a combustible gas for burning in an internal combustion engine coupled to an electrical generator.\nThere are presently many instances where a combustible gaseous fuel, such as methane, is produced as a by-product of other processes, such as in anaerobic digesters. It has long been known that anaerobic digesters can be utilized in decomposing organic waste, such as animal manure, so that the resulting decomposed matter is less offensive and less damaging to the environment. Anaerobic decomposition produces various gaseous by-products including carbon dioxide and combustible methane gas. Until the last decade, because the cost of petroleum and other energy sources has been relatively low, it has not been economically feasible to attempt to collect and utilize the methane produced by the digesters to produce energy such as in electricity. Because of the events of the last decade, particularly concerning the supply of petroleum fuels, the cost of producing energy has increased to an extent where it is now economically feasible to attempt to utilize the methane produced by the anaerobic digesters to produce a usable form of energy. Typically, the methane gas produced by the anaerobic digesters fuels an internal combustion engine which drives a generator for producing electricity. Also, the methane gas can be used as fuel in appliances converted from operation on natural gas.\nA major problem encountered in attempts to utilize the methane gas produced by the digester concerns the need to store the gas when the total amount of gas being produced by the digester is either not used fully by the various appliances utilizing the gas or is not being combusted in the internal gas engine for driving the electrical generator. Prior attempts to overcome this problem have included furnishing a large storage tank with the digester or, at times, including a compressor to compress the gas in volume for easier storage. Another recent example of an attempt to overcome the storage problem is a digester having an expansible cover secured over the decomposing matter. As more gas is produced, the cover expands to a greater volume to store the gas and vice verse. A major problem with this manner of storage aside from the fact that the expansible bag can be punctured and leak is that most users enclose the digester in a suitable building to protect the bag, thereby increasing the expense of the system.\nThe gas storage problem encountered by prior digesters which use methane gas in an engine coupled to a generator for the production of energy is two-fold. First, the associated generator is a synchronous generator which must be rotated at a certain predetermined speed or revolutions per minute (rpm) so that alternating electrical current produced thereby is compatible with utilization by the appliances of the user. Secondly, the consumption of fuel by the generator is determined by the demand placed on it. As a result, the consumption and production rates of gas are different resulting in a need for a storage vessel for any excess gas produced and not contemporaneously consumed by the engine.\nThere are many other types of potential gas supplies which presently are not being utilized but, which also have the attendant problem of fuel storage."}
{"text": "Conventionally, there is known a plasma generation device that generates electromagnetic wave plasma. For example, Japanese Unexamined Patent Application, Publication No. 2010-001827 discloses an ignition device for an internal combustion engine as a plasma generation device of this kind.\nThe ignition device for the internal combustion engine disclosed in Japanese Unexamined Patent Application, Publication No. 2010-001827 emits a microwave generated by a microwave oscillation device to a cylinder, thereby generates low temperature plasma. The low temperature plasma thus generated allows a continuous generation of a large number of OH radicals from moisture in fuel air mixture. In the ignition device for the internal combustion engine, the microwave oscillation device is manufactured as a solid state component."}
{"text": "A lens formed of a high-refractivity low-dispersion glass enables the downsizing of an optical system while correcting a chromatic aberration when combined with a lens formed of a high-refractivity high-dispersion glass. It hence occupies an important place as an optical element for constituting an image-sensing system or a projection optical system such as a projector.\nJP 2007-269584A discloses such a high-refractivity low-dispersion glass. The glass disclosed in JP 2007-269584A has a refractive index nd of 1.75 to 2.00 and has a Ta2O5 content in the range of 0 to 25 mass %, and all of the glasses that have a refractive index nd of at least 1.85 contain a large amount of Ta2O5. That is because the introduction of a large amount of Ta2O5 is indispensable for securing glass stability in the region of high refractivity such as a refractive index nd of 1.75 or more. For such a high-refractivity low-dispersion glass, Ta2O5 is a main and essential component.\nMeanwhile, tantalum (Ta) is an element having a high rarity value and is in itself a very expensive substance. Moreover, rare metal prices are recently soaring worldwide, and the supply of tantalum is deficient. In the field of glass production, tantalum as a raw material is deficient, and if such a situation continues, it may be no longer possible to maintain the stable supply of high-refractivity low-dispersion glasses that are essential and indispensable in the industry of optical apparatuses."}
{"text": "Refuse, recycling and green waste commodity may be placed in any one of a variety of different waste containers. Conventionally, these containers range in volume between two cubic yards and ten cubic yards and can normally be categorized as either being designed for collection by a front end loader (FEL), side loader (SL) or rear end loader (REL) style track. Since the application and methods of collection are considerably different among these styled units, the containers for each of these trucks are normally mutually exclusive.\nFEL waste containers generally include channels or fork pockets built into the sides of the waste container. The collection of commodity from FEL waste containers has typically required a specialized commodity collection vehicle having a pair of spaced forks supported by a pair of lifting arms. To engage and unload the FEL container, the fork must be inserted into the channels and the lift arm must be actuated to lift and invert the FEL container over an opening communicating with an interior storage compartment of the vehicle.\nThe collection of commodity from SL containers requires a specialized collection vehicle having a pocketed assembly mounted to a lift carriage. To engage and lift the SL container, the pocketed assembly is positioned into a specialized hooked plate configuration attached to the SL container and the lift carriage is actuated to lift the SL container over an opening communicating with an interior storage compartment of the vehicle.\nTo load commodity from REL containers requires a specialized collection vehicle having a lower back-end configured to engage each end of a trunnion on the container and a cable at the rear of the vehicle configured to be connected to a top rear portion of the REL container. A motorized winch on the vehicle is then used to pull and lift the REL container off of the ground while the body of the REL container pivots about the trunnion and empties its contents into the rear of the vehicle. In lieu of the winch and cable, some commodity collection vehicles utilize a hydraulic lifting apparatus configured to move the container about the trunnion into a dumping position.\nSince the applications and the methods of collection between FEL, RL and REL tracks are so different, the designs of waste containers to be collected by the different style tracks are normally mutually exclusive. As a result, waste hauling companies are forced to maintain a larger inventory of varying types of collection vehicles and an equally large inventory of varying types of containers making it expensive and inconvenient. Further, the modification of containers to fit the collection mechanisms of other specialized collection vehicles is also expensive and inconvenient. Therefore, there is a continuing economic need to find ways to make these specialized commodity collection vehicles more route flexible by enabling them to be more capable to unload commodity from the normally mutually exclusive containers, wherever possible."}
{"text": "In China lightweight steel structure buildings have been undergo fast development and widely used in industrial buildings. Lightweight steel structure is a new trend for residential buildings. Presently, although the cost is higher than conventional concrete buildings, they are advantageous in their fast construction speed, energy saving and carbon emission reduction. The advantages of lightweight steel structure buildings have been extensively noticed and, they have not only already become the first choice in constructing industrial buildings, but also been extensively used in residential buildings.\nThe structure design, configuration and assembly method of the lightweight steel for buildings have great affects on the quality and construction work. Usually, the main supporting components of the frame/post structure system, such as the structural beams and the structural columns/studs, are made of cold bending or cold rolling light gauge steel having the thickness of 1.5 mm to 5 mm. The cross section of cold bending steel has a uniform thickness and arc-shaped corner. In the bearing wall structure system the cold-formed light gauge steel also comprises various shapes having thickness under 1.5 mm. In addition, structure components are combined with small size of steel angle, steel tube, steel channel, steel plate and steel cylinder, and cold-formed light gauge steel, are also often used as purlin connectors, truss diagonal bracings, etc.\nHowever, there are still a lot to be improved in the building design, structure design, and erection methods of lightweight steel buildings. Presently, when a lightweight steel frame is erected and assembled on site, the supporting components such as the structural beams, the structural columns/studs, and the purlins, etc., made of cold-formed light gauge steel, steel angle, steel channel or steel plates, are usually connected, positioned and fastened by end-to-end joint (such as rigid or hinge). The connecting work is very complicated and serious accumulated errors are brought out due to the large quantity of joint. Therefore, in practical applications, the existing frames for buildings still have lot of defects.\nIn view of the above defects existing in the current lightweight steel structure building, the present inventor creates a novel lightweight steel structure with double continuous beam system by using his rich practical experiences and professional knowledge, combing the application of theory and actively performing the research and innovation, and improves the current lightweight steel structure building thereby bringing more practicality. After iterations of research, experimentation and improvement, the present invention of great practical value is finally created."}
{"text": "1. Field of the Invention\nThis invention relates generally to the field of semiconductor device manufacturing and, more particularly, to a method and apparatus for adaptively scheduling tool maintenance.\n2. Description of the Related Art\nThere is a constant drive within the semiconductor industry to increase the quality, reliability and throughput of integrated circuit devices, e.g., microprocessors, memory devices, and the like. This drive is fueled by consumer demands for higher quality computers and electronic devices that operate more reliably. These demands have resulted in a continual improvement in the manufacture of semiconductor devices, e.g., transistors, as well as in the manufacture of integrated circuit devices incorporating such transistors. Additionally, reducing the defects in the manufacture of the components of a typical transistor also lowers the overall cost per transistor as well as the cost of integrated circuit devices incorporating such transistors.\nGenerally, a set of processing steps is performed on a lot of wafers using a variety of processing tools, including photolithography steppers, etch tools, deposition tools, polishing tools, rapid thermal processing tools, implantation tools, etc. The technologies underlying semiconductor processing tools have attracted increased attention over the last several years, resulting in substantial refinements. However, despite the advances made in this area, many of the processing tools that are currently commercially available suffer certain deficiencies. In particular, such tools often lack advanced process data monitoring capabilities, such as the ability to provide historical parametric data in a user-friendly format, as well as event logging, real-time graphical display of both current processing parameters and the processing parameters of the entire run, and remote, i.e., local site and worldwide, monitoring. These deficiencies can engender nonoptimal control of critical processing parameters, such as throughput, accuracy, stability and repeatability, processing temperatures, mechanical tool parameters, and the like. This variability manifests itself as within-run disparities, run-to-run disparities and tool-to-tool disparities that can propagate into deviations in product quality and performance, whereas an ideal monitoring and diagnostics system for such tools would provide a means of monitoring this variability, as well as providing means for optimizing control of critical parameters.\nOne technique for improving the operation of semiconductor processing line includes using a factory wide control system to automatically control the operation of the various processing tools. The manufacturing tools communicate with a manufacturing framework or a network of processing modules. Each manufacturing tool is generally connected to an equipment interface. The equipment interface is connected to a machine interface which facilitates communications between the manufacturing tool and the manufacturing framework. The machine interface can generally be part of an advanced process control (APC) system. The APC system initiates a control script based upon a manufacturing model, which can be a software program that automatically retrieves the data needed to execute a manufacturing process. Often, semiconductor devices are staged through multiple manufacturing tools for multiple processes, generating data relating to the quality of the processed semiconductor devices.\nVarious tools in the processing line are controlled in accordance with performance models to reduce processing variation. Commonly controlled tools include photolithography steppers, polishing tools, etching tools, and deposition tools. Pre-processing and/or post-processing metrology data is supplied to process controllers for the tools. Operating recipe parameters, such as processing time, are calculated by the process controllers based on the performance model and the metrology information to attempt to achieve post-processing results as close to a target value as possible. Reducing variation in this manner leads to increased throughput, reduced cost, higher device performance, etc., all of which equate to increased profitability.\nCommonly, a processing tool undergoes periodic preventative maintenance procedures or calibrations to keep the tool in optimum operating condition. For example, polishing tools include polishing pads that are periodically conditioned or replaced. Etch tools and deposition tools are periodically cleaned using both in situ cleans or complete disassembly cleans. Steppers are periodically calibrated to maintain alignment accuracy and exposure dose consistency. Many of these preventative maintenance procedures are performed at discrete intervals based on vendor recommendations, past history, and expected degradation rates of consumable items used in the tools. The use of fixed preventative maintenance intervals is not always an effective solution for optimizing tool and line efficiency. If the maintenance activities are performed more often than actually needed, the efficiency of the line and the operation cost of the tool is increased. If maintenance activities are performed less often than needed, product quality and tool reliability may be degraded.\nAnother potential drawback of scheduled maintenance intervals is that unexpected conditions may arise during the time period between maintenance activities. For example, the calibration of a tool (e.g., focus on a stepper) may drift outside expected tolerances. A process controller used to control a particular tool may adjust the operating recipe of the tool trying to correct for what it thinks is normal process variation, while in actuality, the variation is caused by the unexpected condition. Over time, the process controller may adjust its performance model based on post-processing metrology feedback. In some situations, the process controller might not be able to stabilize the process. These control problems may result in increased variation or even defective wafers.\nThe present invention is directed to overcoming, or at least reducing the effects of, one or more of the problems set forth above.\nOne aspect of the present invention is seen in a method for adaptively scheduling tool maintenance. The method includes controlling an operating recipe of a tool using a plurality of control actions, monitoring the control actions to identify a degraded tool condition, and automatically initiating a tool maintenance recommendation in response to identifying the degraded tool condition.\nAnother aspect of the present invention is seen in a manufacturing system including a tool, a process controller, and a tool health monitor. The tool is adapted to process a workpiece in accordance with an operating recipe. The process controller is adapted to control the operating recipe of the tool using a plurality of control actions. The tool health monitor is adapted to monitor the control actions to identify a degraded tool condition and automatically initiate a tool maintenance recommendation in response to identifying the degraded tool condition."}
{"text": "This invention relates to a connector plug for terminating polarization maintaining (PM) optical fibers.\nIn optical fiber communications, connectors for joining fiber segments at their ends, or for connecting optical fiber cables to active or passive devices, are an essential component of virtually any optical fiber system. The connector or connectors, in joining fiber ends, for example, has, as its primary function, the maintenance of the ends in a butting relationship such that the core of one of the fibers is axially aligned with the core of the other fiber so as to maximize light transmissions from one fiber to the other, or, put another way, to reduce insertion loss. Another goal is to minimize back reflections. Alignment of these small diameter fibers is extremely difficult to achieve, which is understandable when it is recognized that the mode field diameter MFR of, for example, a singlemode fiber is approximately nine (9) microns (0.009 mm). The MFR is slightly larger than the core diameter. Good alignment (low insertion loss) of the fiber ends is a function of the transverse offset, angular alignment, the width of the gap (if any) between the fiber ends, and the surface condition of the fiber ends, all of which, in turn, are inherent in the particular connector design. The connector must also provide stability and junction protection and thus it must minimize thermal and mechanical movement effects.\nIn the present day state of the art, there are numerous, different, connector designs in use for achieving low insertion loss and stability. In most of these designs, a pair of ferrules (one in each connector), each containing an optical fiber end, are butted together end to end and light travels across the junction. Zero insertion loss requires that the fibers in the ferrules be exactly aligned, a condition that, given the necessity of manufacturing tolerances and cost considerations, is virtually impossible to achieve, except by fortuitous accident. As a consequence, most connectors are designed to achieve a useful, preferably predictable, degree of alignment, some misalignment being acceptable.\nHowever, in connecting or terminating polarization maintaining (PM) fibers, such is not the case. Many optical fiber components, such as, for example, interferometers and sensors, lasers, and electro-optic modulators, are extremely sensitive to and dependent upon, for proper operation, the polarization of the light. Even very slight alterations or changes in the light polarization orientation can result in wide swings in the accuracy of response of such devices. PM fiber has polarization-dependent refractive indices, and the speed of light in an optical fiber is inversely proportional to the magnitude of the refractive index. A birefringent optical fiber is one having two polarizations having different velocities of propagation, thus giving rise to a xe2x80x9cfastxe2x80x9d wave and a xe2x80x9cslowxe2x80x9d wave. In a PM fiber, the polarization of a linearly polarized light wave input to the fiber, with the direction of polarization parallel to that of the one of the two principal polarizations, will remain or be maintained in that polarization as it propagates along the fiber, hence the term xe2x80x9cpolarization maintaining.xe2x80x9d If the polarization of the light wave is to be maintained at a splice or other connection, the principal axes of birefringence of the two joined fibers must be aligned in parallel, otherwise there will be polarization cross-coupling, i.e., crosstalk, which is highly undesirable. Thus, where two PM fibers, for example, are to be connected together, they should be terminated carefully to reduce the crosstalk during the connectorization process. Also, the connectors must be capable of aligning then maintaining the fiber orientation to the connector key position. Connectors with tolerances adequate for connecting non-PM fibers usually are inadequate for maintaining polarization alignment at the connector junction.\nTypical PM connector requirements are an insertion loss of less than 0.3 dB, and the prior art PM connector arrangements comprise numerous, different connector configurations aimed at meeting these requirements for different connectors, such as an SC type connector as shown in U.S. Pat. No. 5,216,733 of Ryo Nagase et al. The connector of that patent comprises a ferrule body and a ring shaped flange having a keyway mounted on the periphery of the ferrule body. Alignment is achieved by rotating the ferrule body with respect to the flange keyway. The combination of ferrule and flange comprises a plug which is inserted into a push-pull SC connector having a key therein for mating with the flange keyway and springs bias the flange in the longitudinal direction to maintain the alignment.\nIn U.S. Pat. No. 4,784,458 of Horowitz, a splice joint for PM fibers is shown wherein aligned fibers are joined with UV curing epoxy, and the joint is overlaid with epoxy cement for rigidity. Such ajoint is permanent, and does not function as a connect-disconnect optical fiber connector.\nU.S. Pat. No. 5,561,726 of Yao discloses an apparatus for controlling the polarization state of the light within a fiber by squeezing a portion of the fiber to produce a birefringent fiber, and the squeezer is then rotated to change the polarization of the light within the fiber. The device is not a connector, but is intended for use with polarization sensitive devices such as interferometers and electro-optic modulators, however, it may also be used with connectors for connecting two PM fibers.\nIt is common practice in the prior art for creating PM fibers to include a pair of rods in the fiber cladding which extend parallel to the core as shown in U.S. Pat. No. 4,515,436 of Howard et al. Such rods, which are preferably of glass, are, in manufacture of the fiber, included in the fiber preform from which the fiber is drawn. As the fiber is drawn, the rods are accordingly diminished in diameter and are located within the cladding, preferably on either side of the core. The rods have different thermal expansion characteristics than the surrounding glass, and the stress they exert on the core causes the index of refraction to change along that axis. The axes then have different indices of refraction value and thus propagate light at different speeds. Variations on the two rod arrangement are also known, such as the elliptical stress member disclosed in U.S. Pat. No. 5,488,683 of Michal et al. Also, squeezing the fiber to create birefringence, as shown in the aforementioned Yao patent is feasible. The two rod PM fiber, so called xe2x80x9cPandaxe2x80x9d type PM fiber, however, has proven quite satisfactory in use, and it is toward the connectorization of such a fiber that the present invention is directed, although other types of PM fibers may be used with the present invention.\nIn the copending U.S. patent application Ser. Nos. 10/151,450 and 10/151,130 are shown, respectively, an apparatus and methods of tuning a PM connector plug and an adapter for the connector plug of the present invention the principles of which are applicable to any of a large number of optical fiber connectors, but are embodied in a modified LC connector in those applications. For optimum performance, i.e., maximum transmission of a polarized beam, it is highly desirable to provide accurate rotational positioning of better than xc2x11xc2x0 or even as accurate as  less than xc2xcxc2x0 between connectors equipped with polarization maintaining fibers.\nThe present invention is a connector plug for PM connectors which is intended for use with the adapter and tuning method of those Lampert et al. applications to achieve this desideratum. When a PM jumper cable, for example, is terminated by connectors, it is most desirable that the cable/connector combination be tuned to align the fiber slow axis with a reference point such as the connector key which may be the connector latching arm. In accordance with the present invention, there is provided a connector plug which is tunable to yield extremely accurate rotational positioning of the connector ferrules, or, more specifically, the fibers contained therein.\nThe connector plug of the invention comprises in a first embodiment, a tunable barrel assembly as shown in U.S. Pat. No. 5,481,634 of Anderson et al., the disclosure of which is incorporated by reference herein, a connector plug that includes the barrel-ferrule assembly for holding the end of an optical fiber extending axially therethrough and a housing for the assembly. A coil spring member contained within the housing surrounds the barrel, which is of tubular configuration, and bears against an interior wall of the housing and an enlarged flange member on the barrel, thereby supplying forward bias to the barrel-ferrule assembly relative to the housing. As is shown in the aforementioned U.S. patent applications, the barrel-ferrule assembly has an enlarged flange member which is hexagonal 1 in shape and has a tapered or chamfered leading surface that may be slotted. The housing, in turn, has a hexagonally shaped cavity, which provides any of six rotational positions for the flange and a tapered seating surface for the tapered surface of the flange. The dimensions of the cavity are such that the hexagonal barrel flange floats within the hexagonal cavity, in the Anderson et al. arrangement and can rotate about xc2x112xc2x0, which diminishes the tuning accuracy. In the PM connector of the Lampert et al. application Ser. No. 09/811074, this float is greatly reduced by dimensioning the hexagonal barrel flange, or nut, so that it is a virtual slip fit within the hexagonal cavity. Additionally, the flange is affixed to the barrel, hence, the barrel has only six positions, which are subject to the diminished float.\nIn the aforementioned U.S. patent application Ser. No. 09/811074, of Lampert et al, there is also shown a tunable nut arrangement wherein a hexagonal nut is a xe2x80x9clightxe2x80x9d press fit on the ferrule containing tubular portion of the barrel assembly, so that the nut may be rotated relative to the ferrule in extremely small increments, such as fractions of a degree. In the arrangement of that application, the nut may remain seated in the hexagonal seat in the connector plug body, and the ferrule rotated with respect thereto. As explained in that application, the xe2x80x9clightxe2x80x9d press fit is sufficient to prevent accidental or unwanted relative rotation after tuning. After tuning the nut (or flange) may be cemented in place. Such an arrangement is thus capable of producing the small increments of relative rotation that may be required for optimum PM tuning, and is a preferred arrangement for the first embodiment of the invention.\nIn accordance with the present invention, the latch of the plug, which includes a latching arm having latching shoulders on either side, preferably is of plastic material and has tapered sides, thereby imparting to the arm a cross-section shaped as a truncated wedge. The latching shoulders may also have tapered sides so that both the arm and the shoulders, alone or together have a taper. The adapter for use with the invention has a longitudinal slot therein which has sloping sides for receiving the tapered latching arm as shown in application Ser. No. 10/151,130. As is the practice with the latching arm, it is sufficiently resilient or elastic to snap into the slot once the latching lugs are engaged. In prior art arms and slot arrangements the slots and the arm both have straight sides, hence the arm is made slightly narrower than the slot to insure engagement, which gives rise to some play between the plug and the adapter which can alter the position of the polarization vector. The tapered arm and slot of the invention create a tight fit and eliminate such play so that the plug is firmly held in its latched position. In addition, the tapered engagement centers the plug relative to the adapter so that the ferrule/barrel assembly is properly aligned within the adapter.\nFurther in accordance with a second embodiment of the invention, the plug body has a rotatable extender cap therein which has longitudinally extending resilient arms, the distal ends of which bear against the hexagonal portion of the barrel assembly as a three jawed collet further to ensure alignment of the barrel and to eliminate virtually any float thereof. The barrel of the connector plug is an assembly comprising a tubular member which may have a notched end for tuning and an enlarged hexagonal nut affixed thereto on the outer surface of the tubular member and which has a tapered front section which finctions as a guide in conjunction with a tapered recess in the housing. As described hereinbefore, such an assembly is shown and described in the U.S. patent application Ser. No. 09/811074, the disclosure of which is incorporated herein by reference. Tuning of the connector is achieved by rotation of the extender cap which firmly holds the nut on the barrel. Hexagonal seat in the plug body is illuminated, and hence, the ferrule contained in the barrel can be incrementally rotated to produce, where necessary, incremental changes in the angular orientation.\nIn both embodiments of the invention, the latching arm has a tapered cross section as described hereinbefore, and the fitting of the parts to each other approaches a slip fit, so that float is minimized."}
{"text": "Embodiments of the present invention generally relate to stream computing applications. Specifically, the invention relates to the management of processing elements in a streaming data environment, including the predictive start up of processing elements for receiving incoming streaming data.\nWhile computer databases have become extremely sophisticated, the computing demands placed on database systems have also increased at a rapid pace. Database systems are typically configured to separate the process of storing data from accessing, manipulating or using data stored in the database. More specifically, databases use a model where data is first stored, then indexed, and finally queried. However, this model cannot meet the performance requirements of some real-time applications. For example, the rate at which a database system can receive and store incoming data limits how much data can be processed or otherwise evaluated. This, in turn, can limit the ability of database applications to process large amounts of data in real-time."}
{"text": "1. Field of the Invention\nThe present invention relates to a sensor module for detecting an operating parameter of a motor vehicle.\n2. Description of the Related Art\nIt is known to use sensor modules of many different types, in particular in motor vehicles, to acquire various operating parameters, such as pressure sensors, position sensors, air mass sensors, distance sensors, brightness sensors, and the like. The sensor modules convert the physical quantity that is to be acquired, i.e. the relevant operating parameter, into an electrical signal that is supplied to a central or remote data processing device for further processing.\nIn particular, active sensor modules are used that require a power supply, for example in order to acquire the operating parameter or in order to prepare a sensor signal directly at the sensor module in order to enable disturbance-free transmission. In motor vehicles, the supply voltage required for this is transmitted via a supply line that can have a high resistance, for example due to corroded or incompletely joined plug contacts. If, however, an active sensor module is supplied with a voltage that is too low, the signals emitted by the sensor module are faulty or unreliable.\nSensor models are known that recognize undervoltage and overvoltage states of the voltage supply. If the supply voltage falls below or exceeds a fixed, specified voltage threshold, this state is recognized by the sensor module. The sensor module then interrupts the transmission of the sensor signals to the data processing device. Alternatively, the sensor module can send an error signal (instead of the measurement signal). The decision as to whether an undervoltage or overvoltage is present takes place in the sensor module, due to the specified thresholds, and requires a significant circuit outlay. An example of such a sensor module is given in the following.\nPublished German patent application document DE 41 05 283 A1 describes a flow measurement device having a pressure sensor connected to a computing unit. In addition, an undervoltage recognition circuit is connected to this computing unit. If a voltage collapse is detected, the computing unit locally stores all important data. Accordingly, the supply voltage is immediately evaluated locally and protective measures are taken directly at the location in the sensor module. This requires complex circuit parts in the sensor module for the evaluation of the supply voltage and for the taking of the protective measures. In addition, in the sensor module elements must be provided that enable a secure storing. This requires a high-capacitance energy storage device that ensures the supply of energy during the data storing.\nThe complex circuit elements for the voltage evaluation and for the data storage, and of the energy storage device, require complex circuit elements and produce a significant space requirement. Because in particular in automotive technology little space is available for sensors, and the additional circuit elements cause a substantial circuit outlay, an object of the present invention is to reduce the space requirement in comparison with the existing art and to reliably recognize a drop in voltage in a sensor module."}
{"text": "A fuel cell has been proposed as a clean, efficient and environmentally responsible power source for electric vehicles and various other applications. Individual fuel cells can be stacked together in series to form a fuel cell stack for various applications. The fuel cell stack is capable of supplying a quantity of electricity sufficient to power a vehicle. In particular, the fuel cell stack has been identified as a potential alternative for the traditional internal-combustion engine used in modern automobiles.\nOne type of fuel cell is the polymer electrolyte membrane (PEM) fuel cell. The PEM fuel cell includes three basic components: an electrolyte membrane; and a pair of electrodes, including a cathode and an anode. The electrolyte membrane is sandwiched between the electrodes to form a membrane-electrode-assembly (MEA). The MEA is typically disposed between porous diffusion media (DM) such as carbon fiber paper, which facilitates a delivery of reactants such as hydrogen to the anode and oxygen to the cathode. In the electrochemical fuel cell reaction, the hydrogen is catalytically oxidized in the anode to generate free protons and electrons. The protons pass through the electrolyte to the cathode. The electrons from the anode cannot pass through the electrolyte membrane, and are instead directed as an electric current to the cathode through an electrical load such as an electric motor. The protons react with the oxygen and the electrons in the cathode to generate water.\nIt has been desirable to fabricate the fuel cell and related fuel cell components from radiation-cured structures. The formation of radiation-cured structures such as microtruss structures are described in Assignee's co-pending U.S. patent application Ser. No. 12/339,308, the entire disclosure of which is hereby incorporated herein by reference. The formation of radiation-cured fuel cell structures are further described in Assignee's co-pending U.S. patent application Ser. Nos. 12/341,062 and 12/341,105, the entire disclosures of which are hereby incorporated herein by reference.\nRadiation-cured microtruss structures and methodology are described by Jacobsen et al. in “Compression behavior of micro-scale truss structures formed from self-propagating polymer waveguides”, Acta Materialia 55, (2007) 6724-6733, the entire disclosure of which is hereby incorporated herein by reference. One particular method and system of creating radiation-cured structures is disclosed by Jacobsen in U.S. Pat. No. 7,382,959, the entire disclosure of which is hereby incorporated herein by reference. Further radiation-cured structures are disclosed by Jacobsen in U.S. patent application Ser. No. 11/801,908, the entire disclosure of which is hereby incorporated herein by reference.\nTypically, the radiation-cured structures are formed from radiation-sensitive materials such as radiation-curable materials and radiation-dissociable materials. The radiation-cured structure is generally formed by a method including the steps of: providing the radiation-sensitive material; placing a mask between an at least one radiation source and the radiation-sensitive material, the mask having a plurality of substantially radiation-transparent apertures formed therein; and exposing the radiation-sensitive material to a plurality of radiation beams through the radiation-transparent apertures in the mask. The apertures of the mask may be selected to provide different radiation cured structures. To form complex or multi-layered radiation cured structures, different masks having different apertures are generally sequentially applied. Undesirably, masks must be removed following radiation exposure so that further radiation-cured structure may be fabricated. Since masks typically are placed in physical contact with the radiation-sensitive material, the masks must also be cleaned after use in order to remove residual radiation-sensitive material. Although it is possible to achieve multiple radiation exposures through repeated release and application of masks, this methodology can be quite complicated and costly.\nThere is a continuing need for an efficient and cost effective system and method for fabricating radiation-cured structures. Desirably, the system and method provide complex radiation-cured structures without employing costly and inefficient masking techniques, processing steps, and cleaning steps."}
{"text": "The present invention relates generally to medical surgical devices, systems, and methods. More specifically, the invention relates to devices, systems and methods for enhancing knee surgery including measuring forces in flexion and/or extension gap(s) during knee arthroplasty.\nThe knee is generally defined as the point of articulation of the femur with the tibia. Structures that make up the knee include the distal femur, the proximal tibia, the patella, and the soft tissues within and surrounding the knee joint, the soft tissues including the ligaments of the knee. The knee is generally divided into three compartments: medial (the inside part of the knee), lateral (the outside part of the knee), and patellofemoral (the joint between the kneecap and the femur). The medial compartment comprises the medial joint surfaces of the femur, tibia, and the meniscus wedged therebetween. The lateral compartment comprises the lateral joint surfaces of the femur, tibia, and the meniscus wedged therebetween. The patellofemoral compartment comprises the joint between the undersurface of the kneecap or patella and the femur. Four ligaments are especially important in the stability, alignment and functioning of the knee—the anterior cruciate ligament, the posterior cruciate ligament, the medial collateral ligament, and the lateral collateral ligament. In an arthritic knee, protective cartilage at the point of articulation of the femur with the tibia is often worn away, allowing the femur to directly contact the tibia. This bone-on-bone contact can cause significant pain, discomfort, and disability for a patient and will often necessitate knee replacement or knee arthroplasty.\nKnee arthroplasty involves replacing the diseased and painful joint surface of the knee with metal and plastic components shaped to allow natural motion of the knee. Knee replacement may be total or partial. Total knee replacement surgery, also referred to as total knee arthroplasty (“TKA”), involves a total replacement of the distal end of the femur, the proximal end of the tibia, and often the inner surface of the patella with prosthetic parts. Cuts are made on the distal end of the femur and the proximal end of the tibia. Prosthetic parts are then attached. The prosthetic parts create a stable knee joint that moves through a wide range of motion. The replacement of knee structures with prosthetic parts allows the knee to avoid bone-on-bone contact and provides smooth, well-aligned surfaces for joint movement.\nIn TKA, the distal end of the femur and the proximal end of the tibia are completely removed. Thus, healthy tissue may be removed as well. For patients with worn out cartilage in only one compartment, an alternative to TKA may be partial knee replacement, or unicompartmental knee arthroplasty (“UKA”), an example of which is the Oxford® Partial Knee Replacement procedure. Advantages of UKA compared to TKA include smaller incisions, easier post-op rehabilitation, shorter hospital stay, less blood loss, lower risk of infection, stiffness, and blood clots, and easier revision surgery if necessary.\nA challenge in knee replacement surgeries, including both TKA and UKA, is the proper balancing of ligament tension, especially of the medial and lateral collateral ligaments, through the full range of motion of the knee, from a flexed position to a fully extended position. The collateral ligaments, which connect the distal femur and proximal tibia on the medial and lateral sides of the knee, account for much of the stability and movement of the knee. Uneven ligament tension after knee replacement surgery will typically lead to symptoms including joint instability and poor patellar tracking, limited range of motion, impaired function of the knee, and uneven, increased wear and/or loosening of the prosthetic device. These symptoms may often necessitate a repeat surgery. Thus, it is imperative for the short and long-term success of a knee replacement procedure to achieve balanced ligament tension in the knee through the full range of motion of the knee.\nBalancing ligament tension during knee replacement surgery is complicated by the fact that the natural knee does not operate like a hinge moving about a single axis. The knee exhibits dynamic rotation of the tibia relative to the femur as the knee moves from its flexed position to its fully extended position and vice versa. Thus, the natural knee exhibits a rotary laxity that allows the tibia to rotate through a limited internal and external arc during knee flexion. Additionally, the femur translates anteriorly and posteriorly as the tibia is being flexed about it, bringing yet another movement variable. Thus, the ligaments of the knee, along with the femur, tibia and patella, create a dynamic bio-mechanism, making ligament tension balancing in TKA surgery challenging. Many articles and studies have been devoted to ligament tension balancing in TKA, such as: Mihalko, W. H. et al., “Comparison of Ligament-Balancing Techniques During Total Knee Arthroplasty,” Journal of Bone & Joint Surgery, Vol. 85-A, Supplement 4, 2003, 132-135; Eckhoff, D. G. et al., “Three-Dimensional Morphology and Kinematics of the Distal Part of the Femur Viewed in Virtual Reality, Journal of Bone & Joint Surgery, Vol. 85-A, Supplement 4, 2003, 97-104; and Ries, M. D., et al., “Soft-Tissue Balance in Revision Total Knee Arthroplasty,” Journal of Bone & Joint Surgery, Vol. 85-A, Supplement 4, 2003, 38-42.\nOne technique for balancing collateral ligament tension during a knee replacement procedure involves cutting fibers of one or both ligaments to decrease ligament tension—a technique referred to as “ligament release.” A disadvantage of this technique is that the cutting of ligament tissue weakens the ligament(s) and leaves less room for error if future ligament releases or knee replacement procedures are required.\nRather than or in addition to ligament release, the components of a knee prosthesis may be selected and positioned to balance ligament tension. In a typical knee replacement procedure, multiple cuts are made to the distal femur before attaching the femoral component of the prosthesis. Most TKA procedures, for example, involve making distal cuts across the distal end of the femur, anterior and posterior cuts, and angled anterior and posterior chamfer cuts to help secure the femoral component solidly in place. The surgeon attempts to make these femoral bone cuts to achieve a position and orientation of the femoral prosthetic component so as to optimally balance ligament tension through a full range of motion of the knee. However, it is often very challenging to position the femoral bone cuts and femoral prosthetic component to provide ideal ligament tension through the full range of motion of the knee. This is due primarily to the complexity of motion about the knee, as described above, and the difficulty of placing the femoral component so as to maintain desired ligament tension through the full range of motion.\nA number of devices and techniques have been described that attempt to facilitate ligament balancing during a TKA procedure. Some techniques, such as those described in U.S. Pat. No. 5,733,292, involve trial prosthesis components which are used after femoral and tibial bone cuts are made to assess ligament tension. Some devices, such as those described in U.S. Patent Application Publication No. 2003/0187452 A1, are used to measure a gap between the distal femur and proximal tibia in extension and to help a surgeon recreate that same gap when the knee is in flexion. Other “gap checking” devices are described in U.S. Pat. No. 6,575,980. Other devices have been developed to help measure an amount of ligament tension or to apply a desired amount of tension to the ligaments. U.S. Pat. No. 4,501,266, for example, describes a knee distraction device for applying a desired amount of tension. Many paddle-like devices have been suggested for applying or measuring tension across a knee joint, such as the devices described in U.S. Pat. Nos. 5,597,379; 5,540,696; 5,800,438; 5,860,980; 5,911,723; and 6,022,377. Other methods and device are described in co-assigned and co-pending U.S. patent application Ser. No. 10/773,608, Ser. No. 10/973,936, and Ser. No. 11/149,944, the contents of which are fully incorporated herein by reference.\nThe devices described for facilitating ligament balancing during a TKA procedure may also find use in UKA procedures. In certain UKA procedures, for example, the Oxford® Partial Knee Replacement, a surgeon may be provided with a set of shims or feeler gauges, each with a certain thickness. The shims or feeler gauges are used to measure the gap between the distal femur and proximal tibia to facilitate balancing knee tension. Once a side of the distal femur and a complementary side of the proximal tibia has been cut or milled down, the surfaces of the distal femur and proximal tibia may be fitted with template prostheses. With the knee in flexion, the surgeon may measure the thickness of the gap between the femoral and tibial prostheses by placing a shim or feeler gauge into the gap so that “natural tension” in the knee is achieved. “Natural tension” may be achieved when the shim or feeler gauge slides in and out of the gap easily but does not tilt. With the shim or feeler gauge removed, the knee is then extended, for example, in 20 degrees of flexion, and the gap similarly measured. Tension in the knee can be balanced by further cutting or milling off the end of the distal femur so that the gap between the distal femur and proximal tibia is the same when the knee is in flexion as when the knee is extended. The distal femur is further cut or milled based on the difference between the gap size of the knee in flexion and the gap size of the knee extended. For example, if the flexion gap is measured to be 5 mm and the extension gap is measure to be 2 mm, the amount of bone to be milled is 3 mm. This procedure, however, is highly dependent on the subjective “feel” of the surgeon. Thus, there is a need to quantify and standardize the procedure.\nFor the above reasons, a need exists for improved devices, systems and methods for enhancing knee replacement surgery and specifically for dynamically balancing ligaments during knee replacement to improve range of motion, stability, and patellar tracking of the prosthetic knee joint. Ideally, such devices and methods would allow a surgeon to achieve a desired ligament tension balance before committing to and making final bone cuts to the femur. Such devices would ideally be simple to use in conjunction with existing knee replacement procedures and equipment such as prosthesis templates, measurement guides, cutting guides, and saw blades or burs. At least some of these objectives will be met by the present invention."}
{"text": "The present invention relates to optical pitch yaw and roll tracking, and in particular to a system for tracking the position of a head or hand of a user for cursor control in interactive television and computer systems, and also for visual tracking display in virtual reality systems.\nRemote wireless control of a cursor on a television screen or computer screen has become increasingly popular. See, for example, U.S. Pat. No. 5,045,843, incorporated herein by reference. Movements of a photodetector on a movable remote unit are detected with respect to a light beam source on a fixed base unit. Thus, the direction that the remote unit is pointing is determined and this information is used by the base unit to control the position of a cursor on a television screen or computer screen.\nThe photodetector on the movable remote unit detects the pointing orientation with respect to the light beam. This provides for two-dimensional orientation detection and signals related to the angular rotative and/or translational orientation between the light beam source and the photodetector are generated at the output of the photodetector. That is, the \"pitch\" and \"yaw\" of the remote unit are determined with respect to the light beam source of the base unit.\nWhile the photodetector affords adequate pointing position detection of the remote unit relative to the base by determining relative pitch and yaw, it has been found that more accurate pointing orientation detection, and thus more precise cursor control, can be achieved if the \"roll\" of the movable remote unit is determined, to correct the determined \"pitch\" and \"yaw\" components. Moreover, information related to the roll of the remote unit has useful applications alone."}
{"text": "Solution-processed halide perovskites are photoactive materials that may be used to produce low-cost and/or high-performance optoelectronics. For example, perovskite-based solar cells having efficiencies of at least 20%, and perovskite-based light emitting diodes (PeLEDs) having high brightness and tunable color across the entire visible range have been produced (see, e.g., Yang, W. S. et al., Science 2015, 348, 1234; and Tan, Z. K. et al., Nat. Nanotechnol. 2014, 9, 687). PeLEDs typically include organic-inorganic hybrid perovskites, such as CH3NH3PbBr3, as emitters.\nDue to concerns regarding the stability of organic-inorganic hybrid perovskites, all-inorganic halide perovskites, such as cesium lead bromide (CsPbBr3), have been used as alternative emitters in PeLEDs (see, e.g., Brittman, S. et al., MRS Commun. 2015, 5, 7). In addition to a photoluminescence (PL) comparable to that of CH3NH3PbBr3 (Protesescu, L. et al., Nano Lett. 2015, 15, 3692), cesium-based perovskites have exhibited higher thermal and chemical stability (see, e.g., Kulbak, M. et al., J Phys. Chem. Lett. 2015, 6, 2452).\nDespite these advantages, the brightness and/or quantum efficiency of cesium-based PeLEDs remains limited. To improve the external quantum efficiency (EQE), CsPbBr3 quantum dots have been created, which have exhibited higher PL efficiency owing to the surface passivation as well as better morphology in close-packed quantum dot films (see, e.g., Li, G. et al., Nano Lett. 2015, 15, 2640). However, organic capping ligands stabilizing the quantum dots typically suppress charge injection into the emitters, which can result in limited brightness, high turn-on voltage, or a combination thereof.\nAlthough some CsPbBr3 films have been attempted, the relatively low solubility of cesium bromide (CsBr) precursor in commonly used solvents typically makes it difficult to produce CsPbBr3 films of suitable quality. Perovskite-based films that include a large amount of polymer as a host matrix for the perovskite crystallites also have been made, but the inclusion of the host matrix has not cured one or more of the foregoing disadvantages (see, e.g., J. Li, et al., Adv. Mater. 2015, 27, 5196). Due, at least in part, to substantial current leakage caused by incomplete surface coverage, the maximum brightness achieved with PeLEDs based on CsPbBr3 films has been only 407 cd m−2 with peak EQE of 0.008% (Yantara, N. et al., J. Phys. Chem. Lett. 2015, 6, 4360).\nTherefore, methods and materials that overcome one or more of the foregoing disadvantages associated with all-inorganic perovskites are desirable, including methods that [1] reduce or eliminate the substantial current leakage that can occur when films of all-inorganic perovskites, such as CsPbBr3, are used, [2] produce all-inorganic perovskite-based films having favorable physical and/or light emitting characteristics, and/or [3] produce films of all-inorganic perovskites that are stable at a variety of ambient conditions, which may include relatively high humidity."}
{"text": "1. Field of the Invention\nThis invention relates to a refrigerant compressor, and more particularly to a refrigerant compressor for use in an air conditioning system for an automotive vehicle.\n2. Description of the Prior Art\nConventionally, a refrigerant compressor, such as a swash plate compressor, includes a cylinder block formed therethrough with a plurality of cylinder bores, a housing secured to the cylinder block, a swash plate which is tiltably mounted on a drive shaft, for rotation in unison with the drive shaft, a plurality of pistons each reciprocating within a corresponding one of the cylinder bores as the swash plate rotates, and a crankcase defined in the housing. The swash plate is received in the crankcase.\nWhen the compressor is in operation, lubricant collected at a bottom of the crankcase is drawn up and splashed within the crankcase by rotating members including the swash plate as they rotate to be supplied to sliding portions of bearings, the pistons, and other components within the crankcase (so-called splash lubrication method).\nAnother lubricating method of supplying lubricant is a forced-feed lubrication method in which lubricant is forcedly fed to a particular sliding portion by the use of a pump, such as a trochoid pump driven by torque transmitted from a drive shaft of a compressor.\nHowever, when the splash lubrication method is employed, a sufficient amount of lubricant cannot be supplied to unexposed sliding members within the crankcase, such as bearings supporting a rear end of the drive shaft.\nOn the other hand, the forced-feed lubrication method makes it possible to supply a sufficient amount of lubricant even to the unexposed sliding portions within the crankcase. However, this method necessitates provision of a pump, which results in increased manufacturing costs of the compressor. Moreover, it is necessary to secure space for installation of the pump, which inevitably increases the size and weight of the compressor."}
{"text": "1. Field of the Invention\nThe present disclosure is related to mops and mopping tools. More particularly, the present disclosure is related to a mop press for pressing excess liquid from a mop head.\n2. Description of Related Art\nMops are typically used to clean floors or other surfaces by dipping the mop head into a cleaning or rinsing liquid and pressing or wringing excess liquid from the mop head. The damp mop head is then moved across the surface to be cleaned. This process of wetting the mop head, wringing excess liquid from the mop head, and moving the damp mop head across the surface to be cleaned is repeated until complete.\nMop presses for wringing excess liquid from mop heads are known. For example, EP0824008 provides a mop press. Such mop presses typically include a press-plate that is moveable against a stationary wall. In this manner, the excess liquid can be pressed from the mop head by compressing the head between the press-plate and the wall.\nSuch mop presses are typically attached to a bucket or container so that the excess liquid, which is removed from the mop head, is captured in the bucket. Unfortunately, many prior mop presses allow liquid pressed from the mop head to spill outside of the bucket and, thus, back onto the surface being cleaned.\nMany prior mop presses are removably attached to the bucket for ease of use of the bucket. After pressing, the mop head is removed from between the press-plate and the wall. Unfortunately, this action can cause the mop head to catch on the press-plate and/or wall so that the mop press inadvertently becomes unattached from the bucket. Some mop presses have attempted to resolve this problem by providing a locking mechanism for releasably securing the mop press to the bucket. However, such prior art locking mechanisms have been difficult to unlock, which reduces the efficiency of the user of the mop press.\nAccordingly, there is a continuing desire for mop presses that overcome and/or mitigate the above and other deleterious effects of prior art mop presses."}
{"text": "1. Field of the Invention\nThe present invention relates generally to sports equipment, and more particularly to golfing accessories. The present invention is a ball marker adapted for temporary and removable placement on a golf green to mark the location of a temporarily removed ball which might otherwise interfere with the path of another ball in play. The present marker includes various decorative and ornamental features which also provide various functional advantages in durability, visibility, etc.\n2. Description of the Prior Art\nAs the rules governing the game of golf have developed, it has become the practice in most instances to allow a player whose ball lies farthest from the hole on the green, to putt first. Oftentimes, another player's ball may lie on, or very nearly on, the expected path of the farthest ball. Accordingly, the rules of the game provide for the closer player to mark the location of his/her ball temporarily, with a marker that will not impede the travel of another ball if that other ball passes over the marker.\nMost players merely place a dime or other small, flat object (button, etc.) at the location of the temporarily removed ball, to mark its location until it is placed back into play at the proper turn. However, the golfing industry has developed numerous specialized markers for such use, and many players utilize such special purpose markers or some suitable object of particular monetary and/or emotional value to them (rare coin, etc.).\nThese various devices, either developed particularly for the purpose or adapted from other articles, lack various advantages of the present marker invention. One of the greatest problems with such devices is the difficulty in readily seeing and finding the marker when the ball is to be replaced, as generally the present markers are not of particularly high visibility. Moreover, players in a group who may have occasion to place two or more markers down at the same time, may have some difficulty in determining which marker was placed by which player, where common coins or readily available commercial plastic markers are used, as they may all have the same appearance. Also, less expensive markers are often made of metals which may be prone to corrosion, and generally provide no means for attachment to an article of clothing (wrist closure of a glove, lapel buttonhole, etc.). The present marker invention overcomes these deficiencies of the prior art with a marker which is not only of significant value to its owner as a result of the materials incorporated therein, but which also provides superior durability, high visibility, and other advantages over the prior art discussed below.\nU.S. Pat. No. 3,139,690 issued to Miriam L. Hait on Jul. 7, 1964 describes a Golfing Aid comprising a multiple piece magnetic assembly. A garment attachable component includes a pin and magnetic base, and a ball marker component comprises a magnetically attractive disc. While the use of stainless steel is disclosed for corrosion resistance, the device does not provide particularly high visibility, as is provided by the present marker from virtually any angle. Moreover, the Hait marker does not include any depending protrusion to engage the surface of the green, and may thus be easily moved from its resting place by another golfer inadvertently stepping on it or by the travel of a golf ball thereover. The protrusion of the present marker provides reasonable security, and also provides for the attachment of the marker to clothing.\nU.S. Pat. No. 3,895,797 issued to Harrington Moore on Jul. 22, 1975 describes a Golf Ball Marker comprising a small disc with a green engaging peg extending from one side thereof, and a wooly textured fabric material (e.g., Velcro, tm) on the opposite side thereof. The purpose is to provide a texture which reasonably resembles that of the green, in order to avoid any unnatural disruption of the ball should the ball roll over the marker. As a result, the Moore marker is difficult to see, and in fact Moore specifically states that his marker should be unobtrusive in order to avoid distracting the putting golfer. This is unlikely to be a significant distraction, as the golfer is concentrating on the ball and on the hole, rather than on other points, and markers are customarily of relatively small size, e.g., half inch diameter.\nU.S. Pat. No. 3,938,805 issued to Kei Sakuma on Feb. 17, 1976 describes a Golf Accessory comprising a ball marker and attachment means therefor. The attachment means is adhesively secured to another article (the end of a putter, an article of clothing, etc.) and the marker is removably secured thereto by hook and loop material (Velcro, tm). The Sakuma marker is adapted to be placed upon the green with the hook and loop fastening material disposed downwardly, as the marker may include some form of upwardly facing identification means. However, due to the essentially planar configuration of the two Velcro attachment components, no green engaging protrusion is provided by Sakuma. Thus, the Sakuma marker may be easily dislodged from its resting place, as in the case of the Hait marker discussed further above.\nU.S. Pat. No. 4,521,018 issued to N. Eugene Cotchonis on Jun. 4, 1985 describes a Golf Ball Position Marker essentially comprising a C shape, with a central open area slightly larger than the diameter of a golf ball. The Cotchonis marker is adapted for placement over the ball, with the opening passing over the ball, rather than being placed slightly behind and under the rearward part of the ball. Cotchonis relies upon the sheer size of his marker for visibility, rather than providing any highly reflective means for visibility, as provided by the present marker. Also, the Cotchonis marker suffers from the same problem as other markers discussed further above, in that it has no protrusion to engage the surface, and may be moved some distance from its resting place if inadvertently kicked or stepped on by another player.\nU.S. Pat. No. 4,822,052 issued to Gary S. Dimmick et al. on Apr. 18, 1989 describes a Golf Club Grip Attachment For Identification And Ball Position Marking. The attachment is primarily adapted to provide identification for the owner of the clubs, and is preferably permanently affixed to the ends of the club grips. However, the device which is secured to the putter includes a removable ball marker component. The marker component secures removably within the permanently affixed component on the end of the putter by means of a protruding snap fastener-like button in the center of the grip component, and a mating hole in the center of the removable marker component. Thus, the marker does not have any depending point or protrusion adapted to engage the green surface, and may be easily dislodged, as with other similar devices discussed further above. Also, no particularly high visibility is provided for the Dimmick et al. marker.\nU.S. Pat. No. 5,282,616 issued to Marylou I. Stacavich-Notaro on Feb. 1, 1994 describes a Golf Ball Marker closely resembling the Sakuma device discussed further above, in that the '616 marker is removably secured to an underlying component by means of Velcro (tm) or the like. The underlying component may be removably pinned to a garment or the like, or may be permanently secured to the end of a golf putter or other club. The marker itself may include a downwardly extending pin to provide better security when placed on the green. While some identification means is disclosed, no particularly high visibility means is provided in the marker of the '616 patent.\nU.S. Pat. No. D-319,862 issued to Torgny L. O. Lahger on Sep. 10, 1991 describes a design for Green Markers, comprising a disc with a pin protruding centrally therefrom. The pin has a larger diameter protuberance at the distal end thereof. No means of providing high visibility under different lighting angles, as provided by the present invention, is apparent in the Lahger ball markers.\nU.S. Pat. No. D-348,402 issued to Thomas L. McDonald on Jul. 5, 1994 describes a design for a Combined Timepiece And Golfing Tool, comprising a pocket watch type timepiece secured to a double tined turf repair tool. The device would not be used as a ball marker, as it would project vertically from the surface if driven into the green. Even if positioned flat upon the surface, the indicated thickness of the timepiece in the side elevation view would preclude the unimpeded travel of a golf ball thereover.\nBritish Patent Publication No. 664,172 to Jarrett, Rainsford, & Laughton Ltd. et al. and published on Jan. 2, 1952 describes Improvements In Or Relating To Tie Holders, Clips, Or Slides. The clip includes a removable button with an underlying collar adapted to be gripped within the clip. No centrally disposed downwardly extending protrusion adapted to engage a green, or highly reflective surface or other elements, are disclosed in the button element of the '172 British patent. A golfer would not be likely to wear a tie during play, making it unlikely that he or she would have such a clip, with its button, on their person during play. Thus, the button element is not readily adaptable as a ball marker.\nBritish Patent Publication No. 2,049,438 to Masaru Watanabe and published on Dec. 31, 1980 describes a Marker For A Golf Ball, comprising a marker base having a downward extension therefrom adapted to secure to a mating apparel component. The base may have a figure or article therein, with a magnifying lens formed thereover. No particular reflectivity or brilliance is disclosed, as provided by the present marker, nor is any particular durability provided with a relatively soft plastic lens.\nFinally, French Patent Publication No. 2,645,756 to Jean P. Angel and published on Oct. 19, 1990 describes an accessory for golfers comprising a forked turf repair tool including a ball marker removably attachable thereto. The marker has no particular distinguishing features, other than a circular inset element. No particular reflectivity or brilliance is disclosed in the French patent publication, nor is there any particular disclosure of durability of the marker itself (as opposed to the repair tool portion).\nNone of the above inventions and patents, taken either singly or in combination, is seen to describe the instant invention as claimed."}
{"text": "As a human interface for an operation input in the past, a button and a touch panel are the mainstream.\nIn Patent Literature 1, control of an electronic apparatus without the button is described, and Patent Literature 1 discloses that it is not until a proximity sensor reacts that sensing of a touch device with an electrostatic switch starts."}
{"text": "1. Field of the Invention\nThe present invention relates generally to putting devices used in golf, and, more particularly, to a putter a shaft inserted at the rear of a putter head and including a pair of guiding side members.\n2. Description of the Related Art\nAs is well known in the art, a number of different kinds of golf clubs, and especially putters, are known. Relevant to these golf putters, the prior art can generally be grouped into two categories: putters having use-oriented aids and putting stroke training devices.\nPutters having use-oriented aids are disclosed in the art as a means to assist the golfer in using a typical putter, or as an improved putter for use during a game of golf. For instance, an attachment to the shaft of a putter that aids in aligning the golf club at a right angle to a line between the ball and the hole. Others attach to the shaft of a putter and deploy to provide a guide for the putter user. The addition of an indicator surface allows for contrast with the putting surface and helps to avert the problem of the golfer\"\"s eyes attempting to focus simultaneously on the ball, the club face, and the pointer of the putter guide concurrently and thereby losing sight of the pointer device.\nAnother problem that occurs from the use of a typical putter is the turning of the putter head while in use by the golfer. This putting error is caused by the imperceptible turning of the shaft of the putter as the head of the putter approaches or comes into contact with the ball. For example, and illustrating potential problems using a right-handed putter, the slight turning of the putter shaft may result in the toe of the putter passing ahead of the heel of the putter, thereby shutting the putter face and causing the line of the putt to start to the left of the target line. Conversely, slight turning of the putter shaft may result in the heel passing the toe, resulting in an opened putter face and causing the line of the putt to start to the right of the target line. Numerous attempts have been offered to correct this problem.\nOther methods for eliminating putting errors caused by the turning of the putter shaft during use are known. For instance, a putter device that includes an internal spinning gyroscope device inside the putter head to provide sufficient momentum necessary to prevent any xe2x80x9cturningxe2x80x9d of the shaft when the putter head approaches or comes in contact with the ball. However, a putter made in accordance with such references is associated with several drawbacks, the most significant of which is that a putter made in such a manner would appear to not meet all of the requirements and rules of the U.S. Golf Association (USGA), the governing body and policy maker for professional and amateur golf in the United States.\nAdditionally, many putt-training devices are disclosed in the prior art that attempt to assist a golfer in developing proper, consistent habits that generally result from practice conducted with the proper information feedback.\nA search of the prior art did not disclose any patents that read directly on the claims of the instant invention; however, the following references were considered related.\nThe following patents disclose the ornamental design for a golf putting head:\nU.S. Pat. No. D 422,328 issued in the name of Lannoch;\nU.S. Pat. No. D 422,041 issued in the name of Bradford;\nU.S. Pat. No. D 421,471 issued in the name of Moldovan;\nU.S. Pat. No. D 412,729 issued in the name of Woodward;\nU.S. Pat. No. D 410,716 issued in the name of Bendo et al.;\nU.S. Pat. No. D 408,878 issued in the name of Patten;\nU.S. Pat. No. D 399,274 issued in the name of Bradford;\nU.S. Pat. No. D 389,539 issued in the name of Moore;\nU.S. Pat. No. D 385,324 issued in the name of Ubaney; and\nU.S. Pat. No. D 367,909 issued in the name of Lee.\nThe following patents disclose the design and function for a golf putting training head:\nU.S. Pat. No. 5,857,920, issued in the name of Hong, discloses a golf club for putting and driving a golf ball, wherein the club has a reflective surface to properly direct the alignment of a golfer\"\"s eyes to an image of the ball, and also including a leveling device in the head of the club so as to indicate when the putter head is properly aligned in a horizontal orientation;\nU.S. Pat. No. 5,458,332, issued in the name of Fisher, discloses a golf putter head with a cushion face, wherein the face of the putter receives a polyurethane insert or pad at a ball impact position;\nU.S. Pat. No. 5,476,262, issued in the name of Bandiero, discloses a putter trainer having a substantially xe2x80x9cUxe2x80x9d-shaped head with a space provided between two sidewalls so as to receive a golf ball, thereby forcing a golfer to extend the putting follow through approximately five to eight inches in a straight line and providing a simulation of the proper follow through necessary to successfully putt;\nU.S. Pat. No. 4,253,667, issued in the name of Clark et al, discloses a golf ball putter that distributes the mass of the putter head along the line of the putt because the putter shaft is connected to the putter head at a point centrally of the length and width of the head; and\nU.S. Pat. No. 2,472,312, issued in the name of Parrish, discloses a putter directed to impart top spin to a struck golf ball, thereby eliminating the possibility of the golf ball drifting from the intended putting line, wherein the putter has a horizontally elongated head with a pair of legs vertically depending from the ends of the head down to the putting surface.\nConsequently, a need has been felt for providing an apparatus and method of improving putting skills and assisting a golfer in practicing a number of aspects of putting. In particular, a need has been felt for providing a golf putter that accommodates disabled golfers, including golfers that require the use of wheel chairs. Golfers that use wheel chairs often have difficulty properly aligning themselves over the ball, and then making the traditional back-and-through pendulum-type putting stroke used by a majority of golfers. A putter having a rear mounted shaft disposed at an appropriate angle allows one in a wheel chair to position themselves behind the ball, and properly align the putter face and the ball with the golf hole along a more accurate line. Furthermore, the rear mounted shaft allows one to comfortably pull the club back and follow through, thereby averting the potential twisting and turning on the lower back that a standard pendulum putting stroke requires, which might be especially stressful on a person using a wheel chair.\nIt is therefore an object of the present invention to provide an improved golf putting apparatus.\nIt is a feature of the present invention to provide an improved golf putting apparatus including offering an alternative putting stroke, wherein the golfer stands facing the intended line of the putt and placing a back-and-through stroke on the putter.\nIt is a further feature of the present invention to provide an improved golf putting apparatus offering an alternative putting stroke to golfer\"\"s that require the use of a wheel chair, thereby minimizing the twisting and turning of the back.\nBriefly described according to one embodiment of the present invention, a golfing putter is provided comprising a shaft, a putter head and a pair of runner guides. The shaft is a linearly elongated shaft inserted at the rear of the puffer head and disposed at angle to accommodate golfers seated in a wheel chair. The top of the putter head includes an alignment indicia, thereby providing a means for properly aligning the center of the puffer face with the center of the golf ball and on the correct intended putting line. The pair of runner guides are positioned at the bottom of the putter head, intermediate to the putter head and the putting surface, and run parallel to the direction of the shaft and the intended putting line. The rearmost portion of the runner guides includes a rounded surface so as to elevate the putter face, thereby minimizing drag on the putter as it is stroked and impart overspin to the struck golf ball. To use the golfing putter, a golfer is positioned so that the golfer\"\"s dominant arm (or putting arm) is substantially in line with the intended putting line. The user\"\"s arm will grasp the shaft and align the putter face with the intended putting line. The user will then draw the shaft away from the golf ball and then push the shaft into and through the position of the golf ball and toward the golf hole. The putting motion is similar to the striking motions used in games like shuffleboard, billiards, slow-pitch softball or pitching horseshoes.\nAn advantage of the present invention includes reduction and/or elimination of wear and tear inflicted on putting surfaces caused by the need to realign the golfer when using a wheel chair.\nAnother advantage of the present invention is that it is specifically adapted for personal use because of the light weight components and the use of inexpensive materials.\nAnother advantage of the present invention is the reduction of lower back strain for a golfer using a wheel chair, eliminating the necessary reciprocating twists and turns required to perform a pendulum-type putting motion.\nYet another advantage of the present invention is the adaptability of the apparatus for use by persons having a variety of disabilities, including arm or leg amputation, blindness, paraplegia, or musculoskeletal disorders, among others.\nYet another advantage of the present invention is conformity of the present invention to the equipment guidelines established and enforced by the United States Golf Association (USGA), especially within the division that oversees the development and implementation of equipment for disabled golfers."}
{"text": "1. Field of the Invention\nThe present invention relates to tubular connections for application to heavy wall, high performance casing in high pressure, critical service oil and gas wells. More particularly, the present invention relates to such a connection having an improved taper and buttress type threads.\n2. Description of the Prior Art\nIn recent years, oil and gas wells that have been drilled for exploration and production have commonly extended to depths in the range of 15,000 to 35,000 feet, where pressures and temperatures may exceed 15,000 psi and 250.degree. F. respectively. These conditions not only require tubulars (commonly known as Oil Country Tubular Goods (OCTG)) having higher strength, higher reliability and better resistance to failure under higher loads and corrosive applications, such as are disclosed in commonly assigned U.S. Pat. No. 4,354,882, but also require tubular connections which have higher strength and better resistance to severe stress applications. There are three basic types of OCTG's, each of which makes its own contribution to the drilling and completion of oil and gas wells. These types of tubulars are drill pipe, casing and tubing. This invention is primarily intended as a superior connection for heavy wall casing.\nTypically, the lengths of tubulars are threaded at each end and are joined together by an internally threaded coupling. The threads of both the tubes and the couplings must be able to carry the weight of the tubulars and couplings extending below it into the well and must also be able to withstand the high internal and external pressure encountered during drilling and producing of the wells. Since the casing commonly in use will have a weight of from 10 to in excess of 100 pounds per foot, the tensile loads which the threads in the uppermost casing connection must be able to withstand when there are 15,000 to 35,000 feet of tubulars and connections below it are tremendous.\nWith sufficiently high tensile loads, connections may fail due to rupture of the base metal within the connection or by sudden disengagement of the mated male and female threads. The latter of these failures is termed the \"pull-out\" mode of failure.\nThe casing connection must, of course, be leak-proof and a means of sealing must be provided. The threaded connections must be capable of being made up and disassembled without galling and without the danger of stripping threads or cross-threading. It is desirable after the joint is \"stabbed\" that it can be made up with a relatively small number of turns.\nTwo types of casing connections are in common use, coupled and integral joint. Coupled connections employ male threads on each end of the tube which are joined by a female threaded connector called a coupling. Tubes with male threads on one end and female threads on the other are referred to as integral joints since no couplings are required.\nThe following patents and publications illustrate the current state of the art.\nEaton U.S. Pat. No. 2,062,407 discloses a threaded pipe joint of the integral type which employs a buttress type thread having broad crests and roots and which is disclosed to have a reduced galling tendency (page 6, column 2, lines 48-50); easier and faster make-up (page 7, column 2, lines 28-40); and improved hoop strength at the end of the female section of the joint (page 3, column 2, lines 36-46).\nFrame U.S. Pat. No. 2,204,754 discloses a threaded pipe joint of the integral type having a modified acme type thread wherein the thread crests and roots are flat and parallel to the axis of the joint (page 2, column 1, lines 34-42). Five threads per inch are provided and are tapered at both ends, but not in the middle of the thread length, between 3/4\" and 13/4\" per foot (page 2, column 2, lines 21-27).\nRollins U.S. Pat. No. 2,885,225 discloses a drill pipe coupling employing a double thread having three to five threads per inch (column 2, lines 37-43) which is tapered about 11/2\" per foot (column 1, lines 63-67). This taper is said to result in easier and rapid make-up and to promote optimum stabbing (column 1, lines 65-67). The thread is provided with flat crests and roots which are parallel to the taper (column 4, lines 11-20).\nKloesel, Jr., et al., U.S. Pat. No. 3,355,192 discloses a pipe joint of the integral type having a modified buttress type double thread which is tapered 11/2\" per foot (column 3, lines 40-47, and see column 5, lines 40-45).\nYocum U.S. Pat. No. 3,346,278 discloses a tapered V-thread for pipe joints of the integral type wherein the external threads have a basic taper of between 20 percent and 80 percent of the basic taper of the internal threads (column 2, lines 25-28). This arrangement is said to produce a seal between the leading external threads and the innermost internal threads irrespective of whether either thread is at the extreme limits of the commercial taper tolerance (column 3, lines 34-40).\nNowosadko U.S. Pat. No. 3,427,707 discloses a non-threaded soldered pipe joint wherein the female member is machined to a larger taper than the male member so as to assure a seal at the leading end of the male member (column 2, lines 37-41).\nBlose et al. U.S. Pat. No. 3,224,799 discloses a threaded tube joint of the integral type which employs a buttress type thread tapered about 1\" per foot to permit deeper stabbing and more rapid make-up (column 2, lines 51-54; see also column 3, lines 44-47). U.S. Pat. No. 3,224,799 also suggests that the tolerance be such that the actual taper of the male threads will tend to be slightly deeper than that of the female threads (column 3, lines 47-55).\nBlount et al. U.S. Pat. No. 3,307,860 discloses a pipe coupling having an internally thickened central region and shoulders against which a resilient sealing ring is compressed by a shoulder on the ends of the pipe (column 6, lines 36-51; see also column 4, lines 49-55). The threads are of a tapered modified buttress type (column 3, lines 13-15).\nBlose et al. U.S. Pat. No. 3,572,777 discloses a pipe coupling thickened in the central region and provided with an internal groove for a resilient sealing ring which is compressed by the ends of the pipe (column 3, lines 1-18). The thread is of the acme or buttress type and is provided with a relatively steep taper (column 1, lines 5-16).\nAmerican Petroleum Institute (API) Standard 5B (March, 1979) discloses a tubular connection having a tapered buttress casing thread. The API tubular connection consists of a male joining element on the end of the pipe and a coupling member provided with cooperating female joining elements. The threads are on a taper of 3/4\" per foot on the diameter on sizes 133/8\" OD and smaller. The male joining element taper has a tolerance of from +0.0035 inch per inch to -0.0015 inch per inch. The female joining element taper has a tolerance of from +0.0045 inch per inch to -0.0025 inch per inch. The maximum disadvantageous taper spread (i.e., minus on coupling, plus on the pipe end) is 0.006 inch per inch. The API buttress threads have a substantially rectangular form with rounded corners and parallel crests and roots, a 0.200 inch thread pitch, a 3.degree. load flank, 10.degree. trailing flank, and a 0.062 inch thread height with a tolerance of plus or minus 0.001 inch.\nThe thread crests on the API male joining element threads have a length of 0.100 inch with a tolerance of from +0.000 inch to -0.003 inch (+0.000 inch to -0.005 inch on sizes 85/8\" and larger). The rounded corner on the load flank end of the crest has a radius of 0.008 inch with a tolerance of +0.002 to -0.000 and the rounded corner on the trailing flank end of the crest has a radius of 0.030 with a tolerance of from +0.002 to -0.000 inch. The roots have a length of 0.100 inch with a tolerance of from +0.003 to -0.000 inch (+0.005 to -0.000 on sizes 85/8\" or larger). The rounded corners on either end of the root have a radius of 0.008 inch with a tolerance of from +0.000 inch to -0.002 inch.\nThe crests of the API female joining element threads have a length of 0.099 inch with a tolerance of from 0.000 inch to -0.003 inch. The rounded corners at each end of the crest have a radius of 0.008 inch with a tolerance of from +0.002 inch to -0.000 inch. The roots have a length of 0.101 inch with a tolerance of +0.003 inch to -0.000 inch. The rounded corner of the load flank end of the root has a radius of from 0.008 inch with a tolerance of +0.000 inch to -0.002 inch. The rounded corner of the trailing flank end of the root has a radius of 0.030 inch with a tolerance of from +0.000 inch to -0.002 inch.\nThese API buttress thread tolerances result in a thread form that makes contact on the root and crest of the thread and the 3.degree. load bearing flank, but which permits a clearance on the trailing 10.degree. flank of as much as 0.007\" on smaller diameter tubulars. This clearance forms a leakage path that must be sealed by the solids in the thread lube. The API taper of 3/4\" per foot results in 1.984\" of imperfect (non-sealing threads) and, for a 7\", nominal, outside diameter casing corresponds to a buttress thread length of 4.200 inches.\nThe present invention represents an improvement on this API tubular connection and, as compared to that connection, provides a reduced maximum mated thread clearance, combines less imperfect threads with an increase in overall thread length to provide a greatly increased leakage path, and an advantageous taper mismatch which provides high bearing pressure and sealing within the made-up joint adjacent the narrow end of the male member where the coupling is strongest."}
{"text": "1. Field of the Invention\nThis invention relates to a 2-wire control system which communicates with pole mounted push button stations of the kind found at a traffic light controlled intersection via existing pairs of underground wires over which power and data signals are transmitted to enable a visually impaired pedestrian to receive both audible and tactile signals in response to depressing a push button head at a push button station when it is intended for the pedestrian to cross the intersection once vehicular traffic has been halted.\n2. Background Art\nIt has long been known to combine a visual display with a series of traffic lights that are located at an intersection to control vehicular traffic and thereby enable pedestrians to enter the intersection once vehicular traffic has been halted. That is to say, the usual visual display conveys both a written message (i.e., WALK or DON'T WALK) as well as a color sensitive message (i.e., red, green or white) to instruct pedestrians when to cross the intersection. However, such visual warnings are of little value to those pedestrians who are visually impaired. Consequently, a visually impaired pedestrian who activates the push button of a traffic signal will have no way to accurately know when the intersection has been cleared of traffic so that it is time to cross.\nIn order to come into compliance with federal guidelines, such as the Americans With Disabilities Act, cities are required to increase the number of accessible signals that are available to pedestrians at the pole mounted push button stations associated with a traffic light controlled intersection. In particular, to accommodate the needs of visually impaired pedestrians, audible and/or tactile signals are generated at each push button station by which an audible message, a vibration, or the like, is generated when a push button is depressed by one wishing to cross an intersection. In this way, not only will the usual visual message be displayed to sighted pedestrians, but other sensory messages will also become available to coincide with the aforementioned visual message so as to alert visually impaired pedestrians when it is time to cross the intersection after the signal light has changed to halt vehicular traffic.\nIn the past, the pedestrian accessible signaling means has typically been powered at each corner of an intersection by the 115 VAC available at each existing pedestrian lighted sign. Although this approach does not require that additional wires be pulled from each push button station to the traffic control cabinet, the resulting disadvantage is that each push button station operates independently of the others so that sounds cannot be coordinated or synchronized for optimum audible and vibro-tactile presentation to visually impaired pedestrians.\nThe labor costs and the interruption in both vehicular and pedestrian movement at each intersection can be significant as a result of having to install new underground wiring to the push button stations in order to enable the additional signal function to be generated and made accessible to visually impaired pedestrians following the depression of a push button. However, most intersections already contain previously installed pairs of wires that run underground from the existing push button stations to a remote traffic light control box.\nIn this regard, cost sensitive cities would be able to avoid many of the expenditures and inconveniences of having to pull additional wires or even dig trenches and lay new field wires in order to install the new push button stations for each intersection if a control system were available that could incorporate the existing underground wire pairs to transmit power and data signals in order to generate the accessible signal functions for both sighted and visually impaired pedestrians.\nReference may be made to U.S. Pat. No. 5,241,307 issued Aug. 31, 1993 for a microprocessor operated sound signaling and optical signaling generation device that is activated by means of a pedestrian depressing a push button at a traffic light controlled intersection."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an electrical cable connector, and more particularly to a miniature electrical cable connector used in connecting peripheral electronic devices with personal computers.\n2. Description of Prior Art\nFIG. 14 and FIG. 15 show a conventional miniature cable connector comprising an insulative housing 80, a plurality of terminals 90 mounted to the housing 80, a metal shell 70 enclosing the housing 80, a pair of cables connected to the terminals 90 and a cover (not shown) overmolded on the housing 80. The terminals 90 are inserted into the housing 80 from a rear face 801 of the housing 80 through a plurality of passageways 802, respectively. Each terminal 90 has a contacting portion 901 to be connected with a complementary connector, a retaining portion 902 received in a corresponding passageway 802 of the housing 80, and a tail portion 903 resisted by the rear face 801 of the housing 80 to be exposed beyond the housing 80. Thus, a forward movement of the terminals 90 relative to the housing 80 is prevented. However, a hot issue for this kind of connector is that a rearward movement trend of the terminals 90 relative to the housing 80 is inevitable to occur in assembly, especially when soldering conductors of the cable with the tail portions 903 of the terminals 90.\nHence, it is desirable to have an improved structure of the miniature cable connector to overcome the above-mentioned disadvantages of the prior art."}
{"text": "A major problem associated with many electronic circuits is shielding them from unwanted radiation, such as radio frequency interference (RFI) or electromagnetic interference (EMI). This need to shield high-frequency circuits, especially in electronic equipment such as two-way radios, is critical. Conventional shielding systems are generally characterized by a conductive metallic enclosure constructed to completely surround the device or circuit to be shielded. This enclosure acts either to protect the electrical equipment from external RFI or EMI signals or to prevent the escape of RFI or EMI signals generated by the device. Typically, these shielded enclosures are made from a conductive material that is electrically coupled to ground. In prior art, the shielded enclosures have been made by attaching a drawn metallic casing over the affected device(s) and soldering it to a substrate.\nUnfortunately, this method of shielding makes it difficult or impossible to tune or adjust the circuit that is shielded on the interior of the metal enclosure. Circuits containing components that need to be adjusted after assembly, for example, a voltage controlled oscillator (VCO) in a two way radio, cannot be adjusted unless the shield is removed. Removing the shield then affects the circuit and changes the output, thereby rendering any adjustment questionable. Other components, such as laser trimmed resistors, cannot be adjusted without making some type of permanent opening in the shield, again degrading the shielding effectiveness. Most conventional systems try to cover the opening with a separate piece of adhesive-backed metal tape or to solder another piece of metal to the enclosure. Neither of these methods is preferred because they require additional pieces. In the case of the metal tape, the long term integrity of the closure is questionable, and in the case of the soldered patch, another soldering operation is required, exposing the circuitry to additional heat. As a result, a need exists for a method to provide RFI shielding that is economical, reliable, and does not generate excessive temperatures."}
{"text": "Mechanisms used to couple hand tool members such as a head and a handle together are well known. Many of these mechanisms allow members to be coupled and adjusted to a desired position. However, existing devices suffer from many disadvantages. Many of these coupling devices are complex and comprised of many parts. What is desired is a mechanism that is adapted to couple two members together, allows the coupler to be easily engaged and disengaged, and is efficiently made from a minimum number of moving parts."}
{"text": "1. Field of the Invention\nThe disclosure relates generally to a method and apparatus for associating a cable with an electronic device, and more particularly to a method and apparatus for associating a cable with an electronic device and improving electromagnetic compatibility (EMC) shielding between the cable and the electronic device\n2. Description of Background\nAs data rates for cable assemblies increase, better electromagnetic compatibility (EMC) shielding between cable assemblies and electronic devices has become desirable. EMC is the measure of unintentional generation, propagation and reception of electromagnetic energy, and the unwanted effects that such an energy may induce in an electronic component environment. Current methods for improving EMC performance and shielding include cable clamping, cable braid distribution, and use of conductive foams. These methods have offered only marginal improvements, and can be difficult to implement.\nFor example, assembly operations on cables that are to be clamped require that a jacket of the cable be “milked back” in the early stages of cable build. “Milking back” the jacket is a labor-intensive process that requires the operator to manually pull the jacket back thereby exposing a braid. After pulling back the jacket to expose the braid the collapsed area of the jacket is taped to secure it. Subsequently at final assembly of the cable, the jacket is manually pulled back up allowing the back shell to be clamped onto the jacket. This operation is a very difficult, time consuming, and manual process, that is compounded by jacket stiffness, which may vary by bulk cable supplier. With some cable assemblies now requiring plenum rated bulk cable, which includes a very stiff jacket, the current assembly operations requiring jacket “milking” actually cannot be incorporated into assembly. As such, more convenient assembly operations, which eliminate the need for jacket “milking,” and improved EMC shielding, are both desirable."}
{"text": "1. Field of the Invention\nThe invention relates to sliding door and inflatable seal constructions for sterilizing apparatus of the type utilized to sterilize hospital and laboratory glassware, liquids, instruments, parenteral solutions, and paper and fabric materials, such as surgical gowns, sheets and pillow cases, contaminated animal wastes and bedding, and like objects.\n2. State of the Art\nSteam sterilizers are well known and are extensively utilized in hospitals, laboratories and other facilities for sterilizing many types of solid, liquid-containing and porous objects. Typical such sterilizers are illustrated in U.S. Pat. Nos. 3,325,042, 3,339,785, 3,371,986, and 3,488,142.\nThe foregoing sterilizer designs have substantial limitations in their door constructions, in that the door constructions are relatively expensive to make, difficult to operate and do not represent optimum design. For example, in U.S. Pat. No. 3,371,986, the door is indicated as being required to be restrained about three sides and has a reinforced fourth side, in order to maintain an adequate seal and door strength and to prevent the door from being \"blown off\" inadvertently under the substantial pressures of the sterilizer interior. U.S. Pat. No. 3,488,142 is alleged to be an improvement over the design of the foregoing patent, in that it provides restraint of the door about four sides in order to prevent problems that three-sided restraint is inadequate to prevent. The two foregoing designs are, nonetheless, not optimum, in that it is desirable to be able to restrain the door about only two sides, for convenience, ease of fabrication and therefore lower cost, and which is simpler and more convenient to operate.\nIn addition, door seals represent a substantial problem, particularly due to the abrasion of the door seals from glass objects in the sterilizer, some of which may break inadvertently, and glass particles may abrade and even pierce the door seal member, thus creating the inconvenience, \"down time\" and expense of replacement of the door seal member. It is therefore desirable to design a sliding door and seal construction which minimizes the likelihood of door seal wear and replacement, and which facilitates the replacement thereof when necessary.\nU.S. Pat. No. 3,042,980 is the most relevant reference known to applicant from the standpoint of the specific design of the seal construction which is the subject of this patent application."}
{"text": "Three types of opioid receptors, mu, kappa, and delta opioid receptors are generally reported. Recent evidence points to the interactions between receptor dimer combinations of mu, kappa and/or delta receptors (called heterodimers) as also contributing to opioid activity. Opioid receptors and their normal regulation or lack thereof, has been implicated in disease states including irritable bowel syndrome, nausea, vomiting, pruritic dermatoses, depression, smoking and alcohol addiction, sexual dysfunction, stroke and trauma in animals. Therefore it is not surprising that the ability to antagonistically bind opioid receptors has been shown to produce ameliorative, preventative and/or treatment effects in animals including humans afflicted with one or more of these disease states.\nMore recently, certain antagonists of the opioid receptors have been found to increase metabolic energy consumption, and reduction of weight in obese rats while maintaining muscle mass. These findings indicate that an effective opioid antagonist may be useful in preventing, treating and/or ameliorating the effect of obesity. Considering the percentage of the population that is obese in Western societies and the indirect costs associated with treating the effects and symptoms of obesity and Related Diseases, the importance of these findings cannot be overstated.\nThough many opioid antagonists have been disclosed, the search continues for alternative and/or improved or more effective antagonists having an overall benefit to the patient with little or no major side effects. U.S. Pat. No. 4,891,379 disclosed phenylpiperidine opioid antagonists useful for the treatment of diabetes and obesity. In particular, U.S. Pat. No. 4,891,379 disclosed the compound LY 255582 represented by the structure\n\nU.S. Pat. No. 6,140,352 discloses the compound of formula\nFormula 1\n\nwherein the variables X1, X2, X3 R1, R3, R4, R5 and R6 are as described therein, as agonists of the beta adrenergic receptor useful for the treatment of diabetes and obesity.\nRegardless of these and other disclosures of compounds useful as opioid receptor antagonists, or useful for the treatment of obesity, and/or diabetes by other mechanisms, there remains an unmet medical need for a safe, effective and/or alternate treatment or prophylaxis of diseases associated with opioid receptors, particularly obesity and Related Diseases."}
{"text": "Press fit sleeve bearings of both plastic and metal are well known. Metal sleeve bearings are typically pressed into a bore and then machined to the final internal dimension to accommodate the shaft that is intended to operate within the bearing. The post insertion machining is expensive and disrupts the normal flow of assembly of the product.\nThe plastic sleeve bearings have a diffenent drawback. In order for the sleeve to remain in the bore, there must be sufficient deformation in the bearing material to create the forces needed to retain the bearing in the bore. In plastics of the type used for bearings, the required deformation may be quite significant and therefore result in the deformation of the interior dimension. Since it is very diffficult to accurately predict the amount of internal deformation, due to variations in the tolerances of the bore hole and the exterior sleeve dimensions, it is preferable to form the bearing such that the bearing structure permits of controlled deformation. This then minimizes the requirement to perform post assembly operation on the bearing material. It is not very practical to perform a post insertion machining operation on the plastic bearing, since machining of plastic may not be precise and is is difficult to secure smooth bearing surfaces on plastic by machining. The molded skin of a plastic piece provides a superior bearing surface than that of a machined surface.\nPrior attempts to overcome the problem of using plastic materials as bearing sleeves have followed the approach of accommodating the compression of the bearing exterior without unduly effecting the interior bearing surface.\nU.S. Pat. No. 3,359,685 to L. M. Hodgen discloses a bearing sleeve which provides the bearing surface in a series of alternating lands with intermediate grooves, parallel to the axis of the bearing. The lands provide the bearing surface, while the grooves provide flexural relief to accommodate the compression of the material. The exterior of the sleeve is likewise a series of alternating lands and grooves. The lands of the exterior surface are radially aligned to the grooves of the interior surface, thereby permitting the distortion to be accommodated in the region of the groove of the interior surface with a minimum distortion of the bearing surface.\nU.S. Pat. Nos. 1,555,214 to C. W. Johnson and 3,515,417 to J. H. Bowman both use serrations or splines on the exterior surface of the bearing to accept the deformation caused by the insertion of the bearing into the bore of another member. In both of the patents, the splines or serrations extend the full length of the bearing.\nIn the case of the Johnson patent, the bearing is lined with a split sleeve which defines the interior dimension of the bearing. This sleeve may contract or expand. Dimensional stability is thereby sacrificed in the Johnson bearing.\nThe Bowman bearing appears more dimensionally stable on the interior but uses the flexure of the splines to effect a centering of the bearing surface.\nFor precision applications involving relatively light loads, plastic bearings are desirable since they may be molded to very tight tolerances. Cost is very reasonable, particularly if the molded dimensions may be used without further operations in the finally assembled product."}
{"text": "Modern computer systems often deal with large volumes of data. A computer system may have two related processes concurrently executing and processing the set of data. For example, one process may create or retrieve data, while the other process concurrently analyzes or processes the data to accomplish some task. Various traditional algorithms, such as the “ring-buffer” or LIFO stack, are used for buffering the data between the processes. These techniques normally require manipulation of control structures to ensure the integrity of the buffers and chains. These conventional manipulations may include locking of complex structures, spin-loops, or suspensions."}
{"text": "This invention relates to a cold-rolled and annealed steel strip, preferably having a metallic coating thereon to impart corrosion protection to such strip, and exhibiting a high degree of creep resistance.\nThe automotive industry has in recent years sought to improve the performance of automobiles by decreasing the amount of gasoline consumed by automobiles. One aspect of its program of improved performance was to reduce such consumption through a reduction in weight of the automobiles. To compensate for the reduction in weight, by the use of thinner parts, for instance, it was and is necessary to use higher strength materials. By way of example, thinner higher strength low alloy steels are now being substituted for low-carbon, cold-rolled steel. However, since the demands on an automotive component vary due to the component's exposure to high temperatures and/or corrosive conditions, the search for new materials has become very scientific and quite precise. The search for improved materials for automotive exhaust systems represents one of the most challenging needs in the automotive industry.\nThe investigation which led to this invention was undertaken with the goal of developing a sheet steel having improved high temperature strength, and when coated with a metallic coating being resistant to oxidation/corrosion when subjected to the cyclic conditions of an automotive exhaust system at temperatures ranging up to 1500.degree. F.\nTwo steels which have enjoyed some commercial success are Type 409 stainless steel (409SS), and an aluminum coated, titanium-stabilized sheet steel. 409SS, while characterized as a lean stainless steel, i.e. only about 10.5% by wt. chromium, balance essentially iron, it is nevertheless a stainless steel for which a premium is extracted. The titanium-stabilized sheet steel lacked sufficient deformation resistance at elevated temperatures. Even a later innovation on the latter steel by Gupta et al, U.S. Pat. No. 4,398,950, was not sufficient to meet the demands herein stated.\nIt was not until the present invention that titanium, phosphorus and carbon were combined in critical proportions and quantities in steel, and that the role of Ti-P-C concentrations were understood in providing microstructural stability and creep resistance to such steel. This discovery will be described in detail in the specification which follows."}
{"text": "1. Field of the Invention\nThe present invention relates to a machine for automatically assembling parts, and more particularly, to a machine for automatically assembling parts positioned on a carrier in a predefined positional relationship.\n2. Description of the Prior Art\nAn automatic assembling machine conventionally assembles parts onto a chassis supported by a carrier. Generally, one unit supplies the parts and another unit, distinct from the parts supplier, supplies the chassis to the carrier. An automatic assembling machine sequentially picks up the parts from the parts supplier and assembles them on the chassis transported by the carrier.\nAutomatic assembling machines of this type are widely used in manufacturing different kinds of equipment. However, most of these assembling machines are designed and built to assemble only one piece of equipment, and can not be easily changed to accommodate changes in the equipment or changes in the manufacturing process. When a change does occur in the type of equipment to be manufactured or in the manufacturing process, significate changes must be made in the working unit, requiring large expenditures in labor and time.\nThe positioning of the working units with respect to the direction of movement of the moving table which supports the carrier is another drawback of conventional assembling machines. Generally, the working units are not positioned to cooperate with movements of the moving table in the execution of the sequential assembly steps. The movable table must make a number of movements to position the parts with respect to the working units for the subsequent assembly, increasing the time for assembly and contributing to low efficiency of operation."}
{"text": "3-Aminocarboxylic acids (β-amino acids) are important products which are contained in naturally occurring peptides and are important chiral structural units for the synthesis of β-lactam antibiotics, and also for other pharmaceutically active substances, e.g. anti-platelet agents (see, for example, Zablocki et al., J Med. Chem. 36, 1811 (1993), Bovy et al., Bioorg. Med. Chem. 2, 881–895 (1994), Zablocki et al., J. Med. Chem.,38, 2378, U.S. Pat. No. 5,424,334 and U.S. Pat. No. 5,481,021), immunomodulating compounds (see, for example, Suda et al., J. Antibiotics 29, 100 (1976)), hypotensive agents (see, for example, Chaturvedi et al., J. Med. Chem. 13, 177 (1970), Lizuk et al., J. Chem. Soc. Chem. Commun. 1989, 1978, and Okino et al., Tetrahedron Lett. 34, 501 (1993)), or anticancer agents (see, for example, Denis et al., J. Org. Chem. 55, 1957 (1990)).\nThe object of the present invention is the provision of a process which makes possible the preparation of optically active 3-aminocarboxylic acid derivatives in a particularly advantageous manner, in particular in high optical purity and/or high yields."}
{"text": "This invention relates to methods for manufacturing and fabricating low impedance metal gates on MESFETS and more particularly to a method for a fabricating a T-gate for an MESFET employing a double exposure, double develop process.\nIn GaAs MESFETS small gate resistance and short gate length are essential for high gain and low noise performance. The MESFET device is well known and employs a metal semiconductor barrier as a junction. A general approach taken in the prior art for obtaining both low resistance and small length is the adoption of a key-shaped cross section gate (T-gate) structure where the small foot or stem defines the length and the wide top provides a low resistance. The resulting large surface area is essential at operating frequencies above 10 GHZ because of microwave skin effects. T-shaped metal lines have been fabricated by the prior art through the use of a high sensitivity resist on a low sensitivity resist (HI/LO) double layer electronic beam (E-beam) resist system.\nMany such techniques have been investigated in the prior art. As indicated by the prior art, the gate consists of a large trapezoidal cross member fused to a narrow stem and looks like and appears like the letter \"T\". As indicated above, the prior art was well aware of the advantages of such a T-gate and T-gate processing has been employed and described by a number of microwave groups over the past few years. For examples of such prior art, the following references are indicated.\nSee an article entitled \"High Efficiency 35-GHz GaAs MESFETS\" published in the IEEE Transaction on Electron Devices. Vol. ED-34 No. 6 June (1987) by G. C. Taylor, et al. This article describes the fabrication of Ti/Pt/Au gates in a shallow (80 nm) recess etched using gate photoresist as a mask into a 2.5 um long ledge channel. After gate metal deposition, the bottom Ti is chemically etched to a final gate length of 0.4 to 0.5 um. One disadvantage of this technique is that it cannot be used with commonly used Al/Ti gates without losing the top Ti layer. In addition the gate cannot be inspected easily after the Ti etch and ungated channel is guaranteed by the gate etch. Nevertheless the simplicity of the approach and the fact that a ledge is not always required may make the technique feasible if controlled undercut of the Ti layer can be achieved.\nSee also an article entitled \"Electron Beam Fabrication of GaAs Low Noise MESFETs Using a New Trilayer Resist Technique\" published in the IEEE Transaction on Electron Devices, Vol. ED-32 No. 6 June (1985) by P. C. Chao et al. This article describes a LO/HI/LO resist technique utilized to obtain MESFETs with an Ft greater than 100 GHz. The technique requires three resist coatings. A first coating employs a poly (methylmethacrylate) (PMMA). Then a copolymer of methylmethacrylate and methacrylic (MMA-MAA) is employed and finally a thin layer of PMMA is utilized. The total composite thickness is 1 um with a copolymer thickness of 0.5 um and no intermixing between PMMA and the copolymer.\nThe structures described are difficult to obtain given the involved resist chemistry. It is conceded in that paper that the bottom line width is difficult to control. In addition it would appear that even inspection of the patterns prior to recess etching or metalization would also be difficult because of the large aspect ratio between resist thickness and base opening. Thus it is felt that the technique described in that article will not result in an efficient production oriented technique.\nSee an article entitled \"Submicron GaAs Microwave FET's With Low Parasitic Gate and Source Resistances\" published in the IEEE Electron Device Letters, EDL-4 No. 2, February (1983) by S. G. Bandy et al. The technique described here is a straight forward trilayer composite of PMMA/200A Al/PMMA, utilizing only one 20 kV E-beam exposure. The top layer of PMMA is over developed, presumably without effecting the bottom layer, the Al barrier film is etched, and finally the bottom PMMA layer is developed to the required line width. A factor of 6 reduction in gate resistance was obtained leading to devices with a minimum NF of 1.0 db with 13 db of associated gain at 8 GHZ fabricated on n+-n MBE epi in a 0.5 um ledge channel. When utilizing this technique inspection would be difficult although base layer development should be straightforward and the technique should work with Ti/Al as well as Ti/Pt/Au gates. Major drawbacks of the technique is developing variations associated with the porous Al barrier, lack of dimensional control if the top resist chemistry is the same as the bottom, and potential alignment problems with the 200 A Al buried charged dissipation layer.\nSee a further article entitled \"Submicron Lift Off Line With T-Shaped Cross Sectional Form, published in Electronics Letters, June 11, Volume 17, No. 12 (1981) by M. Matsumura et al. This article describes an E-beam, T-gate processing method in which a bi-layer resist composite is used. The top layer is high sensitivity PMIPK (Poly Methyl Isopropanol Ketone) and the bottom layer is a much lower sensitivity PMB. A proprietary developer was employed.\nIn reviewing the article one will ascertain that the technique described is extremely complicated depending critically on well established and reproduceable developing rates. This is one area which has been especially difficult to control in PMMA processing and would quite likely be more difficult to control with a composite resist layer using a single developer. In any event, the technique is not viable to be employed in modern day mass production techniques as they are extremely difficult to implement.\nSee an article entitled \"Synchrotron Lithography for Sub-Half Micron T-Gates in GaAs FET\" published in the Proceedings of Microcircuit Engineering (1986) by K. H. Mueller et al. This technique employed a trilayer PMMA composite consisting of three different molecular weights designed to yield a T-gate stencil after development and exposure to an X-ray dose of 2 to 3 Joules/cm.sup.2. The main problem associated with this technique involved resist processing, in this case optimized for X-ray rather than E-beam lithography. The technique does not appear to be as advanced as the other techniques described in the prior art.\nSee an article entitled \"Fabrication of 0.5 um Plated T-Gates for MESFETs Using the Suss MJB-3 Aligner\", Suss Seminar Series, Publication 101 (1987) by R. W. Baird. This paper describes a very straight forward technique in which a base resist layer is contact printed at 300 nm followed by a recess etching and a sputter deposition of a plating seed (Ti/Pt/Au/Ti). A thicker top resist is then patterned by contact printing to form the top of the T followed by a Ti-etch and gold plating. This technique also proves difficult in that it is dependent on fine line plating techniques and therefore it is incompatible with Al/Ti gates. It is also incompatible with deep recess etching because of the requirement for plating seed continuity.\nSee also an article entitled \"A GaAs Metal-Semiconductor Field-Effect Transistor With a Mushroom Gate Fabricated by Mixed Exposure of Focussed Ion Beams\" published in the Journal of Vac. Sci. Technology, B, Vol. 5 No. 1 January/February 1987 by H. Morimoto et al. on pages 211-214. This article shows a focused ion beam lithography used for fabrication of high resolution patterns. Resist materials with a high resolution such as PMMA which are highly sensitive to ions while showing an extremely low sensitivity to UV light, electron or X-ray are employed. The ion beams are advantageous for high resolution lithography with dimensions less than 0.3 um. These characteristics of the ion beams are suitable for the fabrication of a gate pattern for high performance GaAs microwave FETs. The process flow for the fabrication of such devices is shown in FIG. 1 of article. As one can understand, the above techniques are extremely complicated and are difficult to employ.\nIt is therefore an object of the present invention to provide a T-gate for an FET which includes a deep (0.45 um) recess which increases the height of the stem.\nIt is a further object of the present invention to provide a T-gate for very high frequency power FETs where both short gate lengths and large cross sectional areas are desired.\nIt is a further object of the present invention to provide a method for fabricating a T-gate FET which method employs a thin layer of PMMA for the base of the T-gate and an optical Novolak resist to give a re-entrant side profile performed on top of the T.\nIt is a further object of the present invention to provide a method for forming a T-gate for an FET which method eliminates many of the disadvantages as described in conjunction with the above-noted prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to a data processing method, a reproduction apparatus such as a portable player and an information processing apparatus such as a computer for providing services related to content data.\n2. Description of the Related Art\nFor example, there are systems for managing reproduction history of a plurality of music data (audio data) stored in a hard disk of a personal computer and determining preference of the user based on the reproduction history.\nIn such a system, the user connects a portable player to a personal computer and outputs music data read from a hard disk as the above personal content data to the portable player for storing the same.\nThen, the user reproduces the content data stored in the portable player while the user goes outside.\nHowever, in the above system of the related art, since a user often plays back music data by the portable player, preference of the user cannot be specified correctly only by using reproduction history data in the personal computer, so it suffers from the disadvantage that services according to the users' preferences cannot be provided."}
{"text": "1. Field of the Invention\nThe present invention relates to an intravascular prosthesis, and more particularly, concerns a percutaneously deliverable prosthesis suitable for intravascular reconstruction, reinforcement or repair.\n2. Description of the Prior Art\nBlood vessels of humans or animals undergo a natural degenerative process or are sometimes subject to weakness. In weakened blood vessels, aneurysms may occur. The degenerative effect on blood vessels may also cause a narrowing or constriction of the lumen of the vessel so that blood flow is restricted. In other degenerative situations, or for other reasons, clots or emboli may occur which, should they migrate within the intravascular system, could be very dangerous.\nWith respect to the aforementioned intravascular problems, surgical intervention has been the primary technique for providing relief. For example, aneurysm repair involves a surgical procedure in which an intraluminal vascular prosthesis is inserted into the damaged vessel to reconstruct the section that needs repair. For clogged blood vessels, the excision of thickened atheromatous areas of the vessel has been performed by an endarterectomy procedure. These and other intravascular therapy procedures of an invasive nature are not only risky, but are also costly.\nAngioplasty procedures, using expandable balloons, have been developed for widening the lumen of diseased, constricted blood vessels. Many of these angioplasty procedures are performed percutaneously so that the balloon is introduced into the blood vessel through a catheter inserted through the skin into the vascular system. After the inflation of the expandable balloon widens the clogged blood vessel, it is withdrawn from the blood vessel through the introducer catheter. Balloon catheters are also available in which the inflated balloon is detachable from the catheter once inflated within the blood vessel. The inflated, detached balloon occludes the blood vessel, and is therefore useful in such procedures as varicocele treatment. Such a detachable balloon catheter system, known as the MINIBALLOON.TM. catheter, is sold by Becton, Dickinson and Company, Paramus, N.J.\nAlthough inflatable and detachable balloon catheter procedures, performed percutaneously, are known and available for some intravascular therapy applications, surgery is still relied upon for other applications. Less invasive techniques are being sought for blood vessel repair reconstruction and filtering, as well as vessel occlusion, blood flow regulation or flow assist. The present invention is directed to a device which provides for minimal invasive methods of intravascular therapy. In particular, the present invention is directed to an intravascular prosthesis, percutaneously deliverable which is suitable for the reinforcement or reconstruction of weakened blood vessels, as well as aneurysm repair. Two copending patent applications, having a common assignee herewith, U.S. Ser. Nos. 772,217 now U.S. Pat. No. 4,705,517 and 772,218 both filed on Sept. 3, 1985, relate to those intravascular therapy applications involving blood filtering to prevent clot migration or emboli, and to blood vessel occlusion and blood flow regulation."}
{"text": "1. Field of the Invention\nThe instant invention relates to a sealed diagnostic electrical connector having a retractable seal which retracts to a position to at least partially expose the electrical terminals for probing. The electrical terminals can be probed while the electrical connector is connected to a mating connector, and while the system, to which the connector is applied, is operating.\n2. Description of the Prior Art\nElectrical connectors are often times used outside or otherwise in environments which cause corrosion and the like, which eventually causes discontinuity between at least some of the electrical connections, thus connectors of this type should be sealed, protecting the electrical connection from the hazardous environments and resultant corrosion. Often times when a system becomes inoperable, a possible checkpoint is the electrical connection between the mating electrical connections. However to check the electrical connection, the mating electrical connectors must be disengaged, which could temporarily remedy the situation, as the wiping action between the electrical components during disengagement, could clean the mating electrical components sufficiently enough to temporarily recreate an electrical circuit.\nThus it is desireable to have an electrical connector which is sealed, preventing any stages of corrosion and resultant discontinuity between the electrical mating components. It is also desireable for the electrical connector to include a means of probing the connector, if the system, which is interconnected by the electrical connector, becomes inoperale, without having to disconnect the mating electrical connectors.\nOne method of testing a faulty electrical circuit, presently utilized within the art, is shown in U.S. Pat. No. 3,718,859. This test element is actually a mating connector which is interconnected between the mating electrical connectors. However to utilize such a test element, the mating connectors to which the system is interconnected, must be disconnected, and the test element inserted between the two disconnected mating electrical connectors. If the connection between the mating connectors is faulty, the wiping action on the mating electrical terminals caused by disconnecting the mating connectors may be sufficient to complete the electrical connection between each of the disconnected electrical connectors and between the respective portion of reconnected test element. It may also be possible that an electrical connection can be made between the test element and the electrical connectors, but when the electrical test element is removed and the electrical connectors replaced, the electrical connectors still do not mate properly, or contact each other properly to create an unimpeded electrical connection.\nOther methods of testing inoperable systems include disconnecting the mating connectors which interconnect the system, and testing each half of the system which is disconnected. However, some systems fail to function properly only under certain operating conditions of the system, and therefore to test such a system, the system must be operational while the test is made. If the interconnection between the mating connectors is broken to make the test, the system, is per se inoperable under its normal operating conditions."}
{"text": "1. Field of the Invention\nThis invention relates to apparatus and method for folding-in end closure portions at an end of a packaging sleeve.\n2. Description of Prior Art\nVarious apparatus and methods of folding-in end closure portions of packaging sleeves are known from, for example, U.S. Pat. Nos. 3,820,303; 4,337,059 and 4,524,460; British Patent 1292822; German Patent 687399; Federal German Patent 1952266 and French Patent 7024151.\nFrench Patent 7024151 discloses a system for folding inwards a smaller pair and a larger pair of end closure flaps of a carton. The system includes a compressed-air piston-and-cylinder device which advances, along a supporting rod, an encircling yoke having a smaller pair and a larger pair of arms, the four arms of the yoke being angularly distributed at 90.degree. spacings around the yoke. Mounted diametrically opposite each other upon a ring axially movable along the rod against a spring bias is a first pair of spring-biassed, bell-crank levers, outer free ends of which form a first pair of jaws oscillatable transversely of the axis of the rod. The larger pair of arms acts upon inner free ends of the bell-crank levers to swing the pair of jaws inwards against the spring bias on the levers to fold inwards the smaller end closure flaps. Mounted diametrically opposite each other upon a ring fixed onto the rod is a second pair of spring-biassed, bell-crank levers, outer free ends of which form a second pair of jaws oscillatable transversely of the rod axis. After the larger pair of arms has began to act upon the first pair of bell-crank levers, the smaller pair of arms begins to act upon inner free ends of the second pair of bell-crank levers to swing the second pair of jaws inwards against the spring bias on their levers to fold inwards the larger flaps. Thus, the smaller flaps are folded inwards before the larger flaps. The spring bias on the ring mounting the first pair of levers allows the yoke to advance that ring and the first pair of jaws to remain in their innermost positions while the jaws of the second pair of jaws continue to move towards each other. This system has the disadvantage that the first pair of levers remains in the closure throughout the folding-in of the larger flaps and is thus unsuitable for use in a packaging method in which tucking-tacking of outer sub-panels is required.\nU.S. Pat. No. 3,820,303 discloses a system in which a mandrel wheel transports carton sleeves anti-clockwise through various treatment stations. In a folding station, score lines of bottom end closure panels of the sleeves are \"broken\" by turning inwards a pair of jaws in the form of pivotally mounted, triangular wings, to cause them to bear upon respective inner panels of the end closure. The sleeve is then indexed to a heating station in which selected zones of thermoplastics coating of the bottom end closure are heated to soften the thermoplastics. The heated sleeve is then indexed to a sealing station on the path to which is a tucking device by which the pre-broken panels of the bottom end closure are folded inwards, with a smaller rectangular outer panel thereof becoming tucked under a larger rectangular outer panel thereof. The tucking device includes two folders articulatedly interconnected by a link and turnable, relative to a fixed support plate, about respective axes parallel to the mandrel wheel axis. A first of the two folders is in the form of a U-shaped arm whereof the base of the U is arranged to bear against the smaller outer panel. The second folder is in the form of a bell-crank lever whereof each arm is in the form of a U, the base of the U of one arm being arranged to bear against the larger outer panel and the base of the U of the other arm being attached to a helical tension spring urging this second folder to turn anti-clockwise towards the first folder and thus urging both folders to move oppositely to the advancing bottom end closures. In operation, the second folder is struck by the larger outer panel and thus swung clockwise against the spring bias, thereby, through the link, turning the first folder clockwise to cause it to come to bear on the smaller panel. The design and setting of the tucking device are such that the smaller outer panel is to be folded-in just before the larger panel, so that the smaller panel is tucked under the larger panel. However, during their displacement from the folding station to the tucking device by way of the heating station, the panels can move into relative positions which do not lend themselves to accurate folding and tucking by folders bearing upon only the two outer panels, especially since the carton sleeve is actually moving past the tucking device during the folding and tucking. Moreover, correct initial positioning and return of the folders is dependent upon the tension spring, with the consequence that accuracy of positioning and movement of the folders reduces with increased speed of the mandrel wheel and thus increased production rate.\nU.S. Pat. No. 4,337,059 discloses a tucker-tacker station for folding-in, tucking and tacking the bottom end closures of pairs of carton sleeves when they have been indexed to the station by a rotary turret including pairs of vertical mandrels. The station includes a pair of vertical, cam-operated, reciprocatory, drive rods each of which drives, via a toggle arrangement, a pair of folding jaws arranged to oscillate transversely of the rod axis. The jaws of this pair swing inwards towards each other, transversely of the path of the mandrel pairs, to bear upon and thereby fold inwards the inner pair of panels of the bottom end closure and thus begin the inward folding of the two outer panels of the end closure. Each rod also drives, via a horizontal end plate fixed thereto, a pair of tucker-tacker jaws pivotally mounted upon brackets fixed to the end plate so as to oscillate transversely of the rod axis. The jaws of this pair swing inwards, along a tangent to the mandrel path, to bear on and thus fold inwards the respective panels of the outer pair and thereby to continue the inward folding of the inner pair of panels, one of the latter pair of jaws causing one of the outer pair of panels to lead the other of the outer pair in its inward turning, so that the one becomes tucked beneath the other. At the ends of their inward swings, the tucker-tacker jaws press the panels against the end face of a water-cooled cap of the mandrel to tack them together prior to pressure-sealing at a subsequent bottom sealing station. Because the toggle arrangement passes through dead centre during the lowering of the vertical drive rods and again during the raising of the rods, the pair of folding jaws is moved inwardly and outwardly twice during one bottom end closing operation, with the first inward movement of the folding jaws producing the folding of the inner pair of panels and the second inward movement being superfluous. During the upward movement of the end plate, the tucker-tacker jaws come to bear on the outer pair of panels to fold them inwards, tuck one beneath the other and tack them together. The vertical rods are moved downwards to lower the end plates and then the carton sleeves with tacked bottom end closures are stepped to the pressure-sealing station. This known assembly is relatively complex, with four pivotally mounted jaws."}
{"text": "Certain downhole operations involve placement of elements in a downhole environment, where the element performs its function, and is then removed. For example, elements such as ball/ball seat assemblies and fracture (frac) plugs are downhole elements used to seal off lower zones in a borehole in order to carry out a hydraulic fracturing process (also referred to in the art as “fracking”) to break up different zones of reservoir rock. After the fracking operation, the ball/ball seat or plugs are then removed to allow fluid flow to or from the fractured rock.\nBalls and/or ball seats, and frac plugs, may be formed of a corrodible material so that they need not be physically removed intact from the downhole environment. In this way, when the operation involving the ball/ball seat or frac plug is completed, the ball, ball seat, and/or frac plug corrodes away. Otherwise, the downhole article may have to remain in the hole for a longer period than is necessary for the operation.\nTo facilitate removal, such elements may be formed of a material that reacts with the ambient downhole environment so that they need not be physically removed by, for example, a mechanical operation, but may instead corrode or dissolve under downhole conditions. However, while corrosion rates of, for example, an alloy used to prepare a corrodible article can be controlled by adjusting alloy composition, an alternative way of controlling the corrosion rate of a downhole article is desirable.\nCorrodible materials may include those having a high activity on the saltwater galvanic series, such as a magnesium alloy adjusted for corrosion rate. It has been found that adjusting the amount of trace contaminants in a magnesium alloy can have a significant impact on the corrosion rate of such alloys (Song, G. and Atrens, A., “Understanding Magnesium Corrosion: A Framework for Improved Alloy Performance,” Adv. Eng. Mater. 2003, 5(12) pp. 837-858). For example, metals such as nickel, iron, copper, calcium, etc., may be added to magnesium to increase the corrosion rate and other metals such as zirconium, yttrium, etc. may be added to decrease the corrosion rate. Balancing the amounts of such additives to achieve a desired bulk corrosion rate can in this way control overall corrosion of articles made from the alloy; however, such an approach requires preparation of multiple batches of alloy, requiring high batch-to-batch reproducibility and precise, reproducible control of metal additives or contaminants in the alloy.\nThere accordingly remains a need for controlling the overall corrosion rate of magnesium alloys for use in downhole articles without need for fine adjustment of alloy composition and with improved corrosion control."}
{"text": "The present invention relates to voltage controlled attenuator circuits, and a means to provide accurate gain settings for such attenuators. In particular, the invention relates to such circuits for use in integrated circuit speakerphone systems.\nVoltage controlled attenuators find applications in a wide variety of systems, to adjust the amplitude of signals passed through such systems according to the magnitude of a control signal. For example, in a speakerphone circuit, attenuators are included in both the transmit and receive channels to provide half duplex communication. The transmit and receive attenuators are typically operated in a complementary manner, while one is at maximum gain the other is at maximum attenuation and vice versa. The setting of each attenuator is adjusted so that the difference between the levels remains the same. Using this technique a constant loss is inserted between the two channels and prevents instability that would otherwise occur, due to signal coupling between the loudspeaker and microphone, or sidetone through a hybrid circuit. One requirement of attenuators used in a such an arrangement is accurate gain settings to ensure consistent performance of the speakerphone system. Large variations of the gain level of each attenuator makes the inserted loss between the two channels unpredictable and presents difficulties maintaining the stability of the system.\nIn a typical speakerphone, the gain setting of a channel is dependent upon the detection of speech within that channel. If a far end talker is speaking, the receive signal is greater than the transmit signal, and the transmit attenuator should be set to maximum loss while the receive attenuator is set to maximum gain. The reverse is true if the mirror-end talker is speaking. By monitoring the amplitudes of the signals in both channels, a control circuit may be developed to determine which channel is active and adjust the gains accordingly. A further requirement of an attenuator used in a speakerphone is the need to minimize the feedthrough of control signal into the audio path. Feedthrough generates an audible \"thump\" in the speech channel which occurs when switching from one channel to the other. With sufficient magnitude, the feedthrough causes errors in switching due to its detection as a false speech signal. The primary source of feedthrough is gain dependent offset within the attenuator.\nFIG. 2 illustrates an example of a prior art voltage controlled attenuator. An input voltage is provided at a point 210 through a resistor R2 to a first amplifier 212. Amplifier 212 operates to sink the current from emitter coupled transistors 218, 220 and 214, 216 which receive current from a current source 222. The output is provided through a second amplifier 224 to a voltage output 226. A control voltage indicated VC is applied to transistor 214 and transistor 220. The control voltage VC affects the two transistors in opposite ways, one being NPN, the other PNP. In response to an increasing control voltage, the current in transistor 214 will increase while the current in transistor 220 will decrease. Thus, the current through the two legs of the amplifier section will split in proportion with the control voltage. In accordance with this varying current, the voltage amplification between the input at point 210 and the output at point 226 will vary, such that the gain of the attenuator is directly proportional to the split of the current from current source 222.\nOne problem with a circuit such as that shown in FIG. 2, when used in a typical bipolar integrated circuit, is the different characteristics of the NPN and PNP transistors, which causes a DC feedthrough from the input to the output and which varies in accordance with the level of amplification set by the control voltage. It would be desirable to have a voltage controlled attenuator which eliminates this DC feedthrough effect due to the use of both NPN and PNP transistors in the current splitting core."}
{"text": "Certain motor vehicles, heavy trucks that are powered by diesel engines for example, may be parked for extended periods of time. If the engine is allowed to keep running after a truck has been parked, the alternator should be effective to keep the truck's battery bank charged. If the engine is shut down, either intentionally by the driver, or automatically by a device such as an idle shutdown timer, while the circuits for electric devices remain on, the battery bank will begin to drain through those circuits. For example, an ignition switch will typically remain on after an idle shutdown timer has timed out, and certain circuits that are fed through the ignition switch and may have been left on may continue to draw current.\nBecause the battery bank must be able to provide sufficient current for cranking the engine at starting, the battery bank of a diesel-powered heavy truck typically comprises multiple storage batteries ganged together in parallel. If the battery bank is allowed to drain to the point where it cannot deliver the large cranking current necessary for starting the engine, the vehicle cannot be operated. Because a heavy truck may be parked for an extended period of time, it is therefore appropriate to guard against the possibility that its battery bank will be drained to a point that will not allow the engine to be started or cranked due to failure to turn off all its electric circuits.\nConsequently, it is known to place a disconnect switch in circuit between the battery bank and the rest of the vehicle electrical system for use in disconnecting the battery bank from the electrical system to prevent battery current draw that, if allowed to continue long enough, could drain the battery bank to the point of preventing the engine from being cranked and started.\nCommonly owned U.S. Pat. No. 6,836,094 discloses an example of a battery disconnect circuit that has the advantage of a direct, switch-free connection of the battery bank to the cranking motor so that when the cranking motor solenoid is energized to crank the motor, cranking current flows directly from the battery bank to the cranking motor windings, and not through the disconnect switch.\nBattery disconnect switches are typically electromechanical devices, not solid-state devices. The solid-state counterparts of certain electromechanical devices can often, but not always, provide significant advantages in performance and/or function and/or cost. Sometimes no comparable solid-state counterpart exists for an electromechanical device. For various reasons electromechanical devices continue to enjoy significant usage in the automotive industry, and it is believed reasonable to include battery disconnect switches in that category.\nWhen an electromechanical battery disconnect switch fails, it is typically replaced, but sometimes it is simply bypassed for expediency. The inherent nature of such a switch seems to render it poorly suited for interfacing with vehicle electrical systems that have become largely electronic over the past several decades. Consequently, electronic monitoring of such switches for diagnostic, data collection, and other various purposes appears to have been largely ignored."}
{"text": "The invention relates to an active roll stabilisation system for ships, comprising at least one stabilisation element extending below the water line, which is mounted on a rotary shaft that extends through the ship's hull, sensor means for sensing the ship's movements and delivering control signals on the basis thereof to rotation means for rotating the rotary shaft for the purpose of damping the ship's movements that are being sensed by means of the stabilisation element.\nSuch an active roll stabilisation system for ships is known, for example from U.S. Pat. No. 3,818,959, the disclosure of which is incorporated herein by reference. In said US patent it is proposed to impart a reciprocating rotary motion to a fin-like stabilisation element that projects into the water from the ship's hull below the waterline so as to compensate for the rolling motions that the ship undergoes while sailing. To that end, the ship is fitted with sensor means, for example angle sensors, speed sensors and acceleration sensors, by means of which the angle, the rate of roll or the roll acceleration are sensed. Control signals are generated on the basis of the data being obtained, which signals control the direction of rotation and the speed of rotation of the stabilisation element via the rotation means. Another example of a sensing means is a rate gyroscope, as shown in U.S. Pat. No. 3,756,262, the disclosure of which is incorporated hereby by reference, and also a pendulum-type sensor as shown in U.S. Pat. No. 4,777,899, the disclosure of which is incorporated herein by reference.\nA reaction force acting on the water can be generated by means of said fin-like stabilisation element while sailing, which reaction force imparts a counteracting lifting or torsional moment to the ship, which is to counter the ship's roll, if the stabilisation element is correctly controlled.\nA drawback of the stabilisation system according to said US patent is the fact that it is fairly static as regards the control thereof and that it can only be used while the ship is sailing. The above-described lifting effect does not occur, or not to a sufficient extent, while the ship is stationary, because there is no functional water movement past the stabilisation elements, and consequently there can be no effective roll stabilisation."}
{"text": "There are many situations in which a data stream is transmitted from a transmitter to a receiver over a communication channel. A communication channel may exist in one (or more) of many different possible networks. For example, a communication channel may exist in the internet or in a wireless network such as a mobile telecommunications network. The data in a data stream may be divided into data packets for transmission over a communication channel. The protocol for dividing the data into data packets (e.g. the form of the headers of the data packets, and other such implementation details) may depend upon the nature of the communication channel on which the data is to be transmitted, e.g. the type of network through which the data packet is to be transmitted. The data may be encoded by an encoder prior to transmission from the transmitter and decoded by a decoder after reception at the receiver.\nIn an ideal system, the communication channel is lossless such that every data packet transmitted from the transmitter over the communication channel is received at the receiver. However, in real physical systems, the communication channel may be lossy, meaning that some of the transmitted data packets will be lost on the communication channel and as such will not be received at the receiver. This may be detrimental. Correction schemes can be used at the receiver to help recover at least some of the lost data packets.\nAs an example, forward error correction (FEC) is one correction scheme which can be used to combat packet loss in a communication channel."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device and a method for manufacturing the semiconductor device.\nIn this specification, a “semiconductor device” refers to a device which can function by utilizing semiconductor characteristics; an electro-optical device, a semiconductor circuit, and an electronic device, each of which uses a semiconductor element, are all included in the semiconductor devices.\n2. Description of the Related Art\nIn recent years, high-density integration of a device element has required miniaturization of each element. Here, one of the problems is a contact defect induced by miniaturization.\nA contact defect refers to a malfunction generated when poor connection occurs at a portion where a wiring electrode and a semiconductor device are electrically connected to each other (hereinafter also referred to as contact). In particular, the necessity of electrical connection through a narrow contact hole provided in an interlayer insulating layer has been increased in accordance with a miniaturization technique and multilevel interconnection technique.\nIn the case of a narrow contact hole with a high aspect ratio required by the miniaturization of a semiconductor device, there are problems in that disconnection of a wiring electrode occurs in the contact hole due to a coverage defect of a conductive film for forming a wiring electrode and an upper portion of the contact hole is blocked by the conductive film due to overhang of the conductive film so that the wiring electrode cannot be formed in the contact hole.\nIn order to solve the above problems, as described in Patent Document 1 for example, a method for obtaining electrical connection by using columnar single crystal silicon has been proposed. This method includes the formation of a columnar single crystal of silicon, over a single crystal silicon substrate, by the vapor-liquid-solid (VLS) growth method using a low-melting-point metal material such as gold, followed by doping the columnar single crystal of silicon with an impurity to provide conductivity. The doped columnar single crystal of silicon is used as a connection electrode.\nAs another technique, as described in Patent Document 2, a method has been proposed in which after a contact hole is formed to expose a source region or a drain region of a transistor formed over a single crystal silicon substrate, an epitaxial silicon layer containing an impurity is formed in the contact hole to form a connection electrode."}
{"text": "Processes for removing light components such as hydrogen or nitrogen from a hydrocarbon gas mixture represent a problem in petroleum processing. Typically, the hydrocarbon gas mixture is a refining off-gas or a natural gas and the presence of the light components in the mixture limit further processing, or prevent sales of the gas mixture without the removal of the light components. Often these light components include hydrogen, nitrogen, carbon dioxide, carbon monoxide and mixtures thereof. A wide variety of hydrocarbon gas mixtures are found in petroleum refineries. Some streams are integral parts of specific processes. Often these streams are recycled from a fractionation column or a separator to a reactor. One such recycle stream may be an impure hydrogen stream which must be purified before returning the stream to the reactor. Other of these streams may be by-product streams of a major hydrocarbon conversion process such as a fluid catalytic cracking unit or an ethylene plant. Processes available for the rejection of the light components from these gas mixtures can employ solvent absorption, cryogenic separation, adsorption over molecular sieve adsorbents, or membrane separation. The choice of a suitable process depends upon many factors, some of which are the product purity that is desired, the gas recovery levels, available pressure losses, pretreatment requirements, off-gas composition, impact of reaction products remaining in the light component stream, and the turndown capability of the selected process.\nWhen natural gas is produced from a gas well or is recovered as associated gas from an oil well, the natural gas contains a number of light components which can reduce its heating value. Typically, a natural gas stream comprises nitrogen, methane, ethane, carbon dioxide, inerts, and C.sub.3.sup.+ hydrocarbons. In order to improve the heat content of the product natural gas, the concentration of nitrogen and other inerts must be reduced. This reduction of inerts is often required to meet the quality specifications of pipeline companies that transport natural gas from a well head or natural gas processing plant to the consumer. Typically the natural gas must meet the following specifications:\n______________________________________ Heat Content 900 to 1000 BTU Total Inerts (N.sub.2 + CO.sub.2) 7 mol % maximum Nitrogen 4 mol % maximum ______________________________________\nActual pipeline specifications vary somewhat depending upon the producer's contract for price and quality. In general, a specification for a higher heating value requires a reduced amount of light components such as nitrogen and carbon dioxide. Typically, natural gas at the well head can contain between 3 and 60 mol % nitrogen, between 0.1 and 10 mol % ethane, between 0.1 and 20 mol % C.sub.3.sup.+ hydrocarbons and between 0.1 and 20 mol % CO.sub.2 with the balance being methane. Light components, of which nitrogen is typically the major component, must be removed from natural gas to improve the heat content of the gas and to meet pipeline specification.\nCompositions of the raw gas and the amount of impurities that can be tolerated in the product generally determine the selection of the most suitable process for purification.\nCryogenic technology, which consists of several process variations, has traditionally been employed to separate nitrogen from natural gas. The operating principle of the process entails partially or fully liquefying the high nitrogen content feedstream under pressure and at very low temperatures (e.g., as low as -185.degree. C. (-300.degree. F.)). Afterwards, the partially or fully condensed stream is fractionated in one or two columns, which operate in tandem at two different pressures, to separate the feedstream into a rejected nitrogen stream overhead and a high methane content product stream at the bottom. After heat exchanges with the incoming feedstream, the rejected nitrogen stream is used for gas field operation, vented to atmosphere or reused. The bottom nitrogen-depleted product stream is regasified via process heat exchanges, recompressed, and delivered to the battery limits as an upgraded sales gas.\nU.S. Pat. No. 4,936,888 relates to the use of cryogenic technology in an integrated dual distillation system for rejecting nitrogen in concentrations as high as 75 mol % or more from gaseous methane in a distillation system employing a high pressure fractionator and a low pressure fractionator. The low pressure fractionator accepts a feed predominantly comprising liquid nitrogen under conditions to produce a high purity nitrogen overhead stream and a high purity methane bottom stream.\nU.S. Pat. No. 4,711,651 relates to a process for the separation of a high pressure gas stream such as refining gas in which the starting gas mixture is cooled and separated into a first vapor portion and a first liquid portion which is expanded to an intermediate pressure. The first vapor portion is further cooled and separated into a second vapor portion which may be further processed for ultimate recovery of a methane-rich product gas. Refrigeration is recovered from the mixed intermediate pressure stream.\nU.S. Pat. No. 5,051,120 relates to the cryogenic processing of a feed containing nitrogen and methane. The method is used in treating natural gas which is contaminated with nitrogen. An improved stripping column is provided in which the methane product is produced at a higher pressure than would otherwise be possible.\nOther process options have used non-cryogenic routes which employ solvent extraction. U.S. Pat. No. 4,832,718 relates to a process for processing a natural gas, a thermal or catalytic cracking gas or refining off-gas to produce a methane-rich product, a nitrogen-rich product, a hydrogen-rich stream or an olefin-rich product therefrom by solvent extraction. The process further relates to extractive flashing and extractive stripping using selective physical solvents. The process employs a solvent to separate N.sub.2 from hydrocarbons by extraction (i.e., absorption) at moderate temperatures (ambient to -30.degree. F.) and under pressure. As an N.sub.2 -rich gas stream and a solvent stream come in intimate contact in a packed column, the solvent dissolves nearly all of the hydrocarbons and a small amount of N.sub.2 in the feed gas stream, leaving behind a rejected gaseous N.sub.2 stream that contains the major fraction of the N.sub.2 previously present in the feed stream and a small amount of unrecoverable hydrocarbons. The dissolved hydrocarbons are subsequently released from the rich solvent by successive pressure reductions. The gases that evolve from the pressure reductions are recompressed and recombined into an upgraded sales gas stream. The regenerated lean solvent stream is recycled to the absorber. After refrigeration recoveries, the rejected N.sub.2 stream is available at a relatively high pressure and can be used for field applications, vented to the atmosphere, or reused.\nU.S. Pat. No. 4,623,371 relates to the recovery of hydrocarbons from natural gas containing acidic compounds and from 3-75 mol % nitrogen to provide up to three products: nitrogen-gas product, C.sub.1 -rich gas product, and a C.sub.2.sup.+ liquid product. The process extracts a natural gas stream with a physical solvent to produce a nitrogen stream and a methane-rich solvent stream. The methane-rich solvent stream is flashed to provide a stripped solvent stream and substantially all of the C.sub.1.sup.+ hydrocarbons as a stream of flashed-off-gases. The stripped solvent is recycled to the extraction step.\nU.S. Pat. No. 5,047,074 relates to a process for purging nitrogen from natural gas by passing the gas through an absorption column at elevated pressure and contacting the gas with an absorbent consisting primarily of a poly alpha olefin. The non-absorbed gas consisting predominantly of nitrogen is passed out of the absorption column to waste or to recovery. The resulting rich absorbent is desorbed to obtain the desorbed light hydrocarbons and the lean absorbent.\nU.S. Pat. 5,019,143 relates to the use of a particular group of solvents including paraffins, naphthenes, C.sub.6 -C.sub.10 aromatic compounds and dialkyl ethers of polyalkylene glycol to contact the gas feed mixture in a demethanizing absorber with a reboiler operating between 50 to 400 psig and a temperature of +10.degree. to -40.degree. F. for the separation of ethylene and lighter hydrocarbons from the feed with a distillation column to regenerate the rich solvent. U.S. Pat. Nos. 4,511,381 and 4,526,594 relate to the separation of natural gas from natural gas liquids with physical solvents. Natural gas liquids include hydrocarbons heavier than methane. The physical solvent used for the absorption step in U.S. Pat. No. 4,511,381 is described in terms of a solvent having a relative volatility of methane over ethane of at least 5.0 and a hydrocarbon loading capacity at least 0.25 standard cubic feet of ethane per gallon of solvent. The rich solvent stream is flashed and the gas fraction is compressed, cooled, and condensed to produce the natural gas liquids. U.S. Pat. No. 4,526,594 relates to the selective extraction of a stream of hydrocarbons that are heavier than methane from a natural gas stream with a physical solvent by selectively rejecting the consecutively lowest molecular weight portion of the extracted stream by the use of a second extraction step with a physical solvent to reject hydrocarbons heavier than methane, flashing the resulting rich solvent to separate the selected heavier hydrocarbons, and subsequently de-ethanizing the selected heavier hydrocarbons.\nU.S. Pat. No. 4,883,514 relates to a process for the removal of nitrogen from nitrogen-rich gases which contain more than 3 mol % nitrogen. The nitrogen-rich gas stream is contacted with a lean oil comprised of paraffins, aromatic or cyclohydrocarbons or mixtures thereof having molecular weight between 75 and 250 at temperature no lower than -40.degree. F. to produce a nitrogen stream as an overhead product and a bottoms methane-rich oil stream. The bottoms methane-rich oil stream is flashed to recover a methane-rich overhead gas product and a lean oil rich bottoms stream, and the lean oil stream is recycled to the absorption step. The patent teaches that the contacting or absorption may take place in a methane extraction column which includes a reboiler and is operable in an extractive stripping mode.\nIt is an object of this invention to provide a process for the rejection of nitrogen from natural gas streams using noncryogenic absorption with a hydrocarbon solvent. It is a further object of this invention to provide a process for the rejection of nitrogen from natural gas which can achieve gas recoveries of greater than 99 mol % and provides lower compression requirements, lower capital investment, and lower solvent losses than previous technology."}
{"text": "The present invention relates generally to signal amplifiers, and in particular to signal amplifiers having a non-linear transfer function. The amplifiers of the present invention have a wide range of applicability, including, for example, in mass spectrometry devices or cameras.\nIn many signal processing applications, it is often desirable to digitize analog signals. For example, it is often desirable to convert analog signals representing detected physical events or phenomenon to digital signals for further processing in a digital computer system. To do so, analog signals are typically sampled at a certain rate and converted to digital bits using a digitizer, e.g., an analog-to-digital converter (ADC).\nAn ideal digitization system that is able to handle high dynamic range and high bandwidth signals is difficult to realize. Dynamic range refers to the ratio between the maximum signal level and the minimum signal level that can be handled by the electronic system. Often, the maximum level corresponds to the largest signal that can be handled without clipping or other substantial distortion, and the minimum signal level is determined by the larger of the noise level for small signals or the resolution of the digitizer for small signals. In electrical transmission systems, bandwidth refers to the range between the highest and lowest frequencies of a transmission channel. Bandwidth is typically measured in Hertz (Hz, cycles per second). For a signal sampled at discrete time intervals, the bandwidth of the signal is limited to one half the sampling rate (the Nyquist frequency).\nA mass spectrometer is a good example of a system that requires a digitization system that can handle transient signals from the mass spectrometer that have a high dynamic range and large bandwidth. Ideally the dynamic range of a mass spectrometer designed for protein profiling will match the dynamic range of the concentration of proteins present in blood serum, or about 1015. This dynamic range is many orders of magnitude beyond the capability of any currently available protein profiling system including those using mass spectrometry. The realizable dynamic range of protein profiling systems using mass spectrometry techniques is currently closer to 105. Part of this dynamic range can be practically realized using signal averaging, but much of it must be accounted for directly with the digitization system.\nIn a time-of-flight mass spectrometer (TOF-MS) system, the time resolution of the digitization system is one of the factors that determines the mass resolution that can be realized by the system. Ideally the time resolution of the mass spectrometer is limited by the detection system or the performance of the mass analyzer rather than the digitization system. Currently available ion detectors produce output signals with pulse widths ranging from tenths of nanoseconds to a few nanoseconds. Effectively digitizing these signals without loss of time resolution requires digitization systems with bandwidths from about 100 to about 5000 MHz.\nConsider a digitization system with a dynamic range of approximately 7000 and a bandwidth of approximately 320 MHz. Using standard techniques such a system would require a 13-bit linear analog-to-digital converter (ADC) with a bandwidth greater than 320 MHz. Currently it is difficult and expensive to produce systems with such high combined dynamic range and bandwidth. However, the situation can be improved by considering the noise characteristics of the signal being processed. In general, all analog signals have noise.\nA typical ADC is a substantially linear device that simply outputs a number proportional to the input at a fixed sampling rate. The dynamic range is usually determined by the largest number that the ADC can output. This dynamic range is usually 2n where n is this number of digital bits used to represent the output number. Implicit in this calculation of the dynamic range is the assumption that the smallest signal that can be effectively measured is represented by a change of 1 in the output of the ADC. Thus, for a particular ADC there is a trade off between the largest signal measured and the amplitude resolution achieved. Choice of useful amplitude resolution is generally limited by noise; very little is gained by digitizing a signal with an amplitude resolution finer than the amplitude of the noise on the signal.\nIn general, the noise characteristic of a signal can vary with the amplitude of the signal. Further, the noise characteristic of a signal typically contains contributions from many different sources. In many electronic systems, for example, the amplitude of the noise is essentially independent of the amplitude of the signal (i.e., constant). This is usually true where the main noise contribution is, for example, thermal noise in amplifiers and other circuit components. Another type of noise, shot noise, arises when a signal is composed of discrete elements. For example, any electronic signal associated with a current has shot noise because the current is made up of discrete charge carriers, usually electrons. The amplitude of shot noise is proportional to the square root of the amplitude of the signal. Since most electronic signals involve extremely large numbers of electrons, the shot noise is often a very small fraction of the amplitude of these signals and is often a negligible portion of the total noise in the system. Shot noise also occurs when the signal is due to the detection of particles such as photons or charged or neutral particles. Thus, in a system such as a mass spectrometer, where the numbers of ions detected can be as low as a few or even single ions, shot noise can dominate all other noise sources.\nThe noise characteristic of the signal and the resolution of the digitization system can be matched by choosing the resolution of the digitizer to be roughly equivalent to the noise amplitude (Note that the peak-to-peak noise amplitude will usually be a few times the noise amplitude measured, for example, with a standard deviation). In this way, the largest possible signal is accommodated while the resolution of the system is limited by the noise in the signal and not by the resolution of the digitization (noise limited resolution). This is straight forward for a signal where the noise amplitude is constant. For a linear ADC, this procedure optimizes the dynamic range of the measurement system by minimizing the digitization levels devoted to measuring the noise in the signal and allowing the signal amplitude to be as large as possible. However, for a signal where the noise amplitude changes with signal amplitude, the situation is not so simple. With a linear ADC, matching the resolution to the amplitude of the noise at one signal amplitude will either cause the digitizer to become the resolution limiting factor at signal amplitudes where the noise is smaller or digitizer resolution will be wasted at signal amplitudes where the noise is larger.\nThis situation can be improved by constructing a digitization system where the resolution changes with signal amplitude such that the resolution is always matched to the amplitude of the noise. One way of doing this is to transform the signal before it is digitized with a linear ADC such that the noise amplitude of the transformed signal is a constant, i.e., the noise amplitude does not vary with the amplitude of the signal. For example, for a signal where the noise amplitude characteristic is dominated by shot noise (noise amplitude increases with the square root of the signal amplitude), the required transformation is to take the square root of the signal before digitizing with a linear ADC. Linear and logarithmic transformations are not optimal when the noise characteristic of the signal is dominated by shot noise. For example, with a linear transformation of the input signal before digitization with a linear ADC, the noise amplitude of the transformed signal will increase with larger signal amplitude exactly as it did in the untransformed signal. If the resolution is matched to the noise amplitude at low signal levels where the noise amplitude is smallest, then resolution will be wasted on high amplitude signals where the noise amplitude is larger. With a logarithmic transformation of the input signal before digitization with a linear ADC, the noise amplitude of the transformed signal will decrease with larger signal amplitude. If the resolution is matched to the noise amplitude at high signal levels where the noise amplitude is smallest, then resolution will be wasted on low amplitude signals where the noise amplitude is larger.\nHowever, it is difficult to design amplifiers with particular non-linear transforms, such as a square root transform, and particularly difficult to design such amplifiers to handle high bandwidth signals, for example, signals with bandwidths greater than 100 MHz.\nIt is therefore desirable to provide an amplifier circuit that provides a non-linear transfer function, and particularly a square root transfer function. Such an amplifier circuit should also operate with a high dynamic range and a large bandwidth. It is also desirable to use such an amplifier in a mass spectrometer device or in a camera."}
{"text": "Marketers commonly use databases of customers or potential customers (also referred to as “leads”) to generate personalized communications to promote a product or service. The method of communication can be any addressable medium, e.g., direct mail, e-mail, telemarketing, and the like.\nA marketing database may combine of disparate sources of customer, lead, and/or prospect information so that marketing professionals may act on that information. However, it can be difficult to provide access to a rich set of data in a way that makes sense to the end user of the data (e.g., marketers), as opposed to a database administrator.\nIndeed, despite the existence of current enterprise marketing management suites that purport to integrate marketing databases with campaign management and other marketing tools, it is still common for marketers to hire database consultants to perform common tasks such as defining a market segment and identifying records belonging to the defined market segment. Thus, with existing tools, defining a marketing segment and obtaining a snapshot of its members may be a difficult, expensive, and/or time-consuming process. Moreover, it may be similarly difficult, expensive, and/or time-consuming to update a snapshot of a previously-defined market segment as new records are obtained and/or existing records are modified."}
{"text": "1. Field of the Invention\nThe present invention relates to computer systems and, more particularly, to a system, method and computer program product for managing a computer files system and for providing file sharing among Virtual Execution Environments (VEEs) using file tree systems.\n2. Description of Related Art\nTypically, the operating system of a computer includes a file system to provide users with an interface for working with data on the computer system's disk and to provide the shared use of files to several users and processes. Generally, the term “file system” encompasses the totality of all files on the disk and the sets of data structures used to manage files, such as, for example, file directories, file descriptors, free and used disk space allocation tables, and the like.\nAccordingly, end-users generally regard the computer file system as being composed of files and a number of directories. Each file usually stores data and is associated with a symbolic name. Each directory can contain subdirectories, files or both. The files and directories are typically stored on a disk or similar storage device.\nOperating systems such as UNIX, Linux and Microsoft Windows manage computer file systems by defining a file object hierarchy. A file object hierarchy begins with a root directory and goes down the file tree. The file address is then described as an access path, e.g., a succession of directories and subdirectories leading to the file. This process of assigning a file address is called access path analysis or path traverse. For instance, the path “I/r/a/b/file” contains the root directory (I), subdirectories “r”, “a” and “b” and then the file.\nTypically, the processes within an operating system interact with the file system with a regular set of functions. For example, these functions usually include open, close, write and other system calls or similar mechanisms. For instance, a file can be opened by the open functions and this function acquires the file name as a target.\nThe file system can also include intermediate data structures containing data associated with the file system to facilitate file access. This data is called metadata and may include, for example, data corresponding to the memory location of the files, e.g., where the file is located in the hard drive or other storage medium. For example, in the context of a UNIX operating system, these intermediate data structures are called “inodes,” i.e., index-node. An inode is a data structure that contains information about files in UNIX file systems.\nEach file has an inode and is identified by an inode number (e.g., i-number) in the file system where it resides. The inodes provide important information about files such as user and group ownership, access mode (read, write, execute permissions) and type. The inodes are created when a file system is created. There are a set number of inodes, which corresponds to the maximum number of files the system can hold.\nUsually, computer file systems store this intermediate data concerning the location of stored files as separate structures in the same place as the file content is stored. The functions responsible for file searching, implemented in the operating system kernel, for example, first locate the intermediate data and then locate the file data that is sought. Directories can also have intermediate data structures containing metadata. File systems can also generate intermediate file data “on the fly”, for example, at the moment when the file system is requesting the file. For instance, the NFS Network file system used by Sun Microsystems, provides for on the fly intermediate data creation.\nIn addition, intermediate data structures can include reference files or links that are associated with or point to other files. When a link is accessed, the link itself is not opened. Instead, only the file the link refers to is opened. Thus, the intermediate data structure in a link can contain data referring to other files that are not requested. For instance, the intermediate data structure can contain the path to another file that will be found and opened instead of this reference link.\nThere are several types of links or references. For example, references that include a symbolic name of another file are called “symbolic links.” References that refer to another file's intermediate structure are called “hard links.” The type of link used is generally determined by the operating modes supported by the operating system. File systems can provide several functions. As discussed above, the most basic task of a file system is to provide access to files. File systems can also enhance system performance with additional functions such as, for example, caching, access markers and fault-tolerance.\nThe multi-user operating mode of a computer system can generally allow the operating system processes of different users to operate simultaneously. Each process within the operating system is usually associated with information that identifies the user. For instance, in a UNIX system, this information is typically an identifier of the user and group on whose behalf this process is being executed. When accessing a file, the operating system defines the user requesting the file operation and determines whether the operation is permitted for that user.\nGenerally, this determination can be made upon opening the file, e.g., using a file access mechanism of the OS, such as requesting a function of the type “open.” Thus, on the basis of this access information, the operating system can organize different views of the same file system tree based upon selected parameters, such as, for example, time, operation type or user information.\nTo unite different types of computer file systems, these file systems can be mounted. For any directory inside the file system, it is possible to mount another file system into that existing directory. Thus, one tree of the computer file system appears inside another file tree. The operating system uses a specific system call of the operating system kernel to mount a file system. This system call includes at least two arguments: the mounting point, e.g., the directory inside of the current file system, and the file system itself, e.g., the storage device or memory location where the data resides. Depending on the file system, additional information containing parameters of the specific file system types can be included.\nDuring analysis of the access path to the selected data file, the operating system defines a moment when the path “passes” through this mounting point and “below” this point, an analysis is performed using operations specific for the given file system. The set of these operations is defined according to the parameters established during the file mounting process.\nThe UnionFS file system, developed in the FreeBSD UNIX operating system, implements a similar technique. One feature of UnionFS is that each user can have a different view of the tree of the same file. In order to provide this feature, two trees of the file system are built when mounting UnionFS. The first tree is a read-only tree. The second tree is created during the user's session and is used for auxiliary purposes. This second tree is defined as an additional parameter when mounting.\nWhen calling a file within the shareable tree, a search is performed in two ways. First, the search can be based on a path name that is computed based on the location of the file. For example, the mounting point of UnionFS can be located at “a/b/u,” and the file to be addressed can be at “/a/b/u/c/d/e.” The second tree, mounted to the same point, is located starting from the address “/x/y/.” Then an additional address is computed as “/a/b/u/c/d/e” minus “/a/b/u” plus “/x/y/.” As a result, the additional address is computed as “/x/y/c/d/e.”\nThus, the specific intermediate data structure (e.g., inode) is searched using the computed path name. If the specific intermediate data structure (inode) is found, then it is assumed that the file is found and the requested operation will be performed on this file. If the file is not found, then a second search by the direct address is provided. If the file is not found there either, the system returns the corresponding error code. Otherwise, the system acts according to the requested operation. If the file opens in response to an operation to modify its content or associated data, then the file is first copied to the computed address as described above, and the operation is performed on the new copy. Otherwise, the operation is performed on the file located in the shareable tree by the requested address.\nOne way to change the search address of the file object, and, accordingly, the position of the root file system for a group of processes, is to use a primitive that is analogous to the OS UNIX kernel primitive “chroot.” The operation of this primitive is based on the principle of shifting the real root of the file system or “root” directory to some location for a selected group of processes, for instance, for all processes of one user. Then all file operations inside this process kernel are performed only within the sub-tree of the selected file system.\nAnother example of this type of system is one based upon “snaps” of the file system, or tree snapshots, in which modifications to the entire file system are chronologically arranged. All modifications made in the file system or any of its parts during a period of time are saved in a separate tree of the file system. Such separate chronologically arranged trees represent the complete history of file system modifications for a discrete period of time. Thus, to determine the file state at a fixed moment of time, the operator searches for the file in the most recently accessed file tree. If the file is not located, then the previous tree is searched.\nSimilarly, the Mirage File System (MFS) from IBM, describes a system consisting of a number of trees and a specific file search mechanism that depends on the file type, extension and sequence of requests, among other parameters. One of the principles of this computer file system is the substitution of the file search path whereby the search path is expanded to other file system locations associated with the file object being searched. For example, this system offers an implementation of a system of snapshots.\nU.S. Pat. No. 6,289,356 also describes an example of implementation of specific intermediate data structures, in which the file system is organized with a strictly regulated mode of modifications records. The disclosed system provides the transition of file system states so that at any moment of time the system is in the correct state. Additionally, the system creates snapshots of the file system through doubling an intermediate data structure (e.g., inode) without doubling the files themselves. The system also marks the files chosen to store data file blocks as belonging to some snapshot of the file system. This provides interference with file system functioning at the level of data distribution algorithms.\nThe industry trend of virtualization and isolation of computer system resources makes the task of file sharing more complex. A Virtual Execution Environment (VEE), such as, for example, Virtual Machine (VM) or Virtual Environment (VE) (for example, a Virtual Private Server (VPS)), is a type of an isolated VEEs that run on the same physical machine simultaneously. Each VEE instance accesses its own files, as well as (in some cases) files used by other VEEs.\nVirtualization allows running a number of VEEs (such as VMs or VEs/Virtual Private Servers) on the same physical machine or processor. Thus, sharing file data among numerous VEEs executing on the same machine becomes even more crucial. A robust file sharing system is especially important in multi-user systems, such as, for example, virtual server systems. A VE is a server, for example, a Web server, that shares computer resources with other virtual servers.\nInstead of requiring a separate computer for each server, dozens of virtual Web servers can co-reside on the same computer system. In most cases, performance is not affected and each web site behaves as if it is being served by a dedicated server. It is also desirable to provide VEs with an access to shareable files residing on network storage.\nIn addition to maintaining efficient allocation of resources, providing multi-user access in Virtual Execution Environment involves other considerations as well, including security, avoiding file corruption and maximizing system efficiency. Accordingly, it is desirable to have a method for providing file sharing among Virtual Execution Environments (VEEs) using file tree system that provides multi-user access to files residing on the local storage as well as on the remote network storage."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent the work is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure.\nReferring now to FIG. 1, a functional block diagram of a wireless communications system is presented. A home network 102 receives packets from and sends packets to a packet data network (PDN) 104, such as the Internet. A wireless terminal 106 wirelessly connects to the home network 102. For example, the wireless terminal 106 may be a mobile phone, and the home network 102 may be the cellular network of a mobile phone operator.\nThe wireless terminal 106 is configured to work with the home network 102 and may be unable to connect to the networks of other carriers. In various implementations, the wireless terminal 106 may be able to view content from the PDN 104 via the home network 102. The home network 102 may interconnect with the networks of other service providers.\nReferring now to FIG. 2, a functional block diagram of a wireless communications system offering mobility is presented. The home network 102 is connected to one or more visited networks 110-1, 110-2, . . . , and 110-X (referred to herein as visited networks 110). In various implementations, the visited networks 110 may be the networks of other service providers, including service providers in other countries.\nIn FIG. 2, the mobile wireless terminal 120 has established a wireless connection to the visited network 110-1. The mobile wireless terminal 120 includes the code and data used to communicate with the home network 102 via the visited network 110-1. In this way, the mobile wireless terminal 120 can interface with the home network 102 even when connected to one of the visited networks 110.\nThe mobile wireless terminal 120 includes mobility features that allow the mobile wireless terminal 120 to communicate with the visited networks 110. The mobile wireless terminal 120 may use a host-based Internet protocol (IP) when handling mobility management. A mobility protocol may include a mobile IP (MIP), which may refer to a host-based IP or a network-based IP. Dual stack mobile IP (DSMIP) is an exemplary host-based mobility protocol. Several versions are available. For example, DSMIPv6 includes a mobile wireless terminal 120 that extends its IP stack when changing its point of attachment. The mobile wireless terminal 120, according to DSMIPv6 is also involved in signalling that enables IP mobility. DSMIPv6 is often referred as client mobile IP (CMIP). DSMIPv6 is described in request for comment (RFC) 3775, titled “Mobility Support in IPv6,” the disclosure of which is hereby incorporated by reference in its entirety.\nReferring now to FIGS. 3A-3B, a functional block diagram and timeline of an implementation of host-based mobility are presented, respectively. The mobile wireless terminal 120 roaming in a visited network 110-1 connects through a home agent (HA) 180 of the home network 102 to a packet data network (PDN). The home agent 180 may be included in a gateway (GW) and establishes a logical location of the mobile wireless terminal 120. Packets ultimately destined for the mobile wireless terminal 120 are sent to the home agent 180, while packets from the mobile wireless terminal 120 will appear to originate from the home agent 180.\nFIG. 3B shows an example of how the mobile wireless terminal 120 manages its mobility by using DSMIPv6. For example, with DSMIPv6, the mobile wireless terminal 120 is secured using a security protocol, such as Internet key exchange (IKEv2), which is the protocol used to set up a security association (SA) in the Internet Protocol Security (IPsec) suite. IPsec includes a suite of protocols for securing IP communications by authenticating and/or encrypting each IP packet in a data stream.\nThe mobile wireless terminal 120 may register a current location of the mobile wireless terminal 120 with the home network 102 and thus receive session continuity during roaming by using DSMIPv6. To register the current location, the mobile wireless terminal 120 may require a home address, a home agent address and a security association with the home agent 180. The procedure to obtain such information is referred to as “bootstrapping”, and the mobile wireless terminal 120 may use various information protocols, such as Directory Name Service (DNS) or Dynamic Host Configuration Protocol (DHCP), to obtain such information dynamically during bootstrapping.\nWhen the mobile wireless terminal 120 attaches to a visited network 110-1, the mobile wireless terminal 120 first performs layer 2 and layer 3 attach procedures. Layer 2 and layer 3 attach procedures are procedures of the Open Systems Interconnection Basic Reference Model (OSI Model). During the layer 2 and layer 3 attach procedures, the mobile wireless terminal 120 acquires a topology-correct IP address, commonly referred to as a care-of address (CoA). For the layer 2 attach procedure, the mobile wireless terminal 120 may establish a connection to an access router (AR) 182 within the visited network 110-1. In various implementations, additional access routers (not shown) may be present. The access router 182 may communicate with other networks, including the home network 102 and/or a corporate network (CN).\nWhen the mobile wireless terminal 120 connects to the visited network 110-1, the mobile wireless terminal 120 performs authentication and authorization with the access router 182. This may include communicating with an authentication, authorization, and accounting (AAA) server in the home network 102. The AAA server may retrieve information based on an identifier of the mobile wireless terminal 120, such as a network address identifier, that uniquely identifies the mobile wireless terminal 120. The AAA server may indicate to the access router 182 whether the mobile wireless terminal 120 should be allowed access to the home network 102 and what services should be provided to the mobile wireless terminal 120.\nUpon authentication, the mobile wireless terminal 120 receives a local address from the access router 182. Using this local address, the mobile wireless terminal 120 may implement layer 3 attach procedures to communicate with various network elements, including the home agent 180. The mobile wireless terminal 120 may then initiate a bootstrapping procedure by providing an access point name (APN) during either DHCP or DNS to obtain, for example, a home agent address. An access point name may include a logical name that identifies a Packet Data Network (PDN). The mobile wireless terminal 120 may include a list of access point names or other identifiers for various PDNs.\nThe mobile wireless terminal 120 may then establish a security association with the home agent 180 through, for example, IPSec IKEv2, in order to obtain a home address (HoA) from the home network 102. The mobile wireless terminal 120 may send the information of the CoA and HoA in a binding update message to the home agent 180. Upon receiving the binding update message, the home agent 180 may “bind” the CoA and the HoA in a binding cache. In one implementation, the home agent 180 creates a binding cache entry that records information about the mobile wireless terminal 120, including the current address of the mobile wireless terminal 120.\nThe home agent 180 may send a binding acknowledgement (Ack) message to notify the mobile wireless terminal 120 of the status of the binding update. The binding acknowledgement message may also include the HoA so that the mobile wireless terminal 120 is aware of a global home address assigned to the mobile wireless terminal 120. The mobile wireless terminal 120 may therefore “connect” to the home agent 180 via the exchange of the binding update and binding acknowledgement messages. The home agent 180 may use the binding cache entry to forward packets from destinations associated with the home network 102 to a current point of attachment of the mobile wireless terminal 120.\nFor example, packets sent from other network elements, such as other mobile wireless terminals, are sent to that HoA. The home agent 180 receives those packets and forwards the packets to the mobile wireless terminal 120. Similarly, packets from the mobile wireless terminal 120 are first forwarded to the home agent 180. The home agent 180 then forwards the packets with the source address of HoA. To allow for packets to be exchanged between the mobile wireless terminal 120 and the home agent 180, a tunnel may be established between the mobile wireless terminal 120 and the home agent 180.\nReferring now to FIGS. 4A-4B, the mobile wireless terminal 120 may set up the connectivity to multiple PDNs 104-1, 104-2 (collectively referred to as PDNs 104) via multiple home agents in respective PDN GWs 316-1, 316-2 (collectively referred to as PDN GWs 316) using DSMIP. The timeline of FIG. 4B includes the mobile wireless terminal 120 attaching to each of the PDN GWs 316 via a procedure similar to the one shown in FIG. 3B. In other words, the mobile wireless terminal 120 uses separate DSMIP bootstrapping procedures to connect to each of the PDN GWs 316.\nEach PDN GW 316 may include one or more home agents that each have respective addresses. However, for simplicity, the terms PDN GW and home agent may be used interchangeably herein. The mobile wireless terminal 120 may establish tunnels to either or both the home agents to route traffic to the appropriate home agent based on the service associated with that traffic. For example, one PDN 104-1 may be used to provide a service such as push e-mail, while another PDN 104-2 may be used to provide voice over internet protocol (VoIP) telephony."}
{"text": "The Quinolin-2(1H)-one derivative is a heterocyclic compound having a specific substituent or substituents is known as an active ingredient that has serotonin uptake inhibitory activity (or serotonin re-uptake inhibitory activity) in addition to dopamine D2 receptor partial agonistic activity (D2 receptor partial agonistic activity), serotonin 5-HT2A receptor antagonistic activity (5-HT2A receptor antagonistic activity), and adrenaline α1 receptor antagonistic activity (α1 receptor antagonistic activity). This active ingredient has a wide therapeutic spectrum for central nervous system diseases (particularly schizophrenia). In the pharmaceutical field, the development of pharmaceutical preparations that are suitable according to the severity of disease in various patients in need of treatment, patient predisposition, and other factors, has been desired.\nBrexpiprazole, which is a Quinolin-2(1H)-one derivative, having the chemical name 7-{4-[4-(1-Benzothiophen-4-yl)piperazin-1-yl]butoxy}quinolin-2(1H)-one a compound of formula I has been approved under the trade name Rexulti® as an Tablet, oral having the following dosage forms 0.25 MG, 0.5 MG, 1 MG, 2 MG, 3 MG, 4 MG.\n\nU.S. Pat. No. 7,888,362 B2 (U.S. Pat. No. '362) discloses Brexpiprazole and its pharmaceutically acceptable salts. U.S. Pat. No. '362 discloses the following scheme for the preparation of Brexpiprazole.\nwherein X, Lg represents a halogen atom or a group, which causes a substitution reaction the same as in halogen atom.\nThe key intermediates in the preparation of Brexpiprazole are 7-Hydroxyquinoline-2(1H)-one a compound of formula II and 1-(benzo[b]thiophen-4-yl)piperazine hydrochloride of formula IV. One of the intermediate in the preparation of Brexpiprazole, 7-Hydroxyquinoline-2(1H)-one a compound of formula II was synthesized from Journal of Medicinal Chemistry Volume 58 Issue 14 Pages 5561-5578, WO 2014085284 A1, WO 2012003418 A2, WO 2002002108 A1 wherein discloses the preparation of Formula II by the reaction of N-(3-methoxyphenyl)cinnamamide with aluminum chloride at reflux temperature.\nThe provided prior art processes involves higher temperatures for the cyclization reaction as well as involves class 2 solvents, such as chlorobenzene, which have limitations in the use of these solvents as per the FDA ICH guidelines.\nCN 104844585 A discloses the following scheme for the preparation of Brexpiprazole:\n\nCN 104447723 A discloses the following scheme for the preparation of Brexpiprazole:\n\nThe present inventors have followed the above provided process and observed that if DDQ used in the final stages of the reaction impurities are formed, which are difficult to be removed and in turn results in low yield and low purity of the obtained product.\nAnother key intermediate involved in the preparation of Brexpiprazole is 1-(benzo[b]thiophen-4-yl)piperazine hydrochloride and was synthesized as follows:\nwherein X is a leaving group.\nHowever the above reaction involves in the formation of dimer of formula X as impurity during the condensation of compound of formula XI and XII.\n\nU.S. Pat. No. '362, WO 2007026959 A2, WO 2008047883 A1, JP 4540700 B2, JP 4785881 B1, Organic Process Research & Development, Volume: 19, Issue: 4, Pages: 555-558, CN 104529998 A, WO 2015054976 A1, CN 104829602 A, CN104892589 A discloses the reaction of 1-(benzo[b]thiophen-4-yl)piperazine with conc. hydrochloric acid, wherein the purity of the obtained 1-(benzo[b]thiophen-4-yl)piperazine hydrochloride salt is low with significant amount of dimer impurity, which inevitably decreases the yield of final compound i.e., Brexpiprazole.\nFurther the prior art processes discloses the column purification methods to increase the purity of Brexpiprazole, however it is not commercially and industrially feasible.\nU.S. Pat. No. 9,206,169 B2 (U.S. Pat. No. '169) discloses the preparation of 1-(benzo[b]thiophen-4-yl)piperazine by reacting compound of formula XI and XII in the presence of (a) a palladium compound and a tertiary phosphine or (b) a palladium carbene complex, in an inert solvent or without a solvent in order to suppress the by-product impurity. However, the process is not economical and cost-effective as the used palladium complexes are not easily available. Further after the reaction, the removal of by-products is not easy.\nUS 20150087655 A1 discloses Brexpiprazole dihydrate, anhydride and its process for the preparation thereof.\nCN 104844586 A discloses Brexpiprazole amorphous form and its process for the preparation thereof.\nCN 104829603 A discloses Brexpiprazole hydrochloride monohydrate and process for the preparation, wherein designated as Form A.\nPolymorphism is defined as “the ability of a substance to exist as two or more crystalline phases that have different arrangement and/or conformations of the molecules in the crystal Lattice. Thus, in the strict sense, polymorphs are different crystalline structures of the same pure substance in which the molecules have different arrangements and/or different configurations of the molecules”. Different polymorphs may differ in their physical properties such as melting point, solubility, X-ray diffraction patterns, etc. Although those differences disappear once the compound is dissolved, they can appreciably influence pharmaceutically relevant properties of the solid form, such as handling properties, dissolution rate and stability. Such properties can significantly influence the processing, shelf life, and commercial acceptance of a polymorph. It is therefore important to investigate all solid forms of a drug, including all polymorphic forms, and to determine the stability, dissolution and flow properties of each polymorphic form. Polymorphic forms of a compound can be distinguished in the laboratory by analytical methods such as X-ray diffraction (XRD), Differential Scanning calorimetry (DSC) and Infrared spectrometry (IR).\nBrexpiprazole and its hydrochloride salt can exist in different polymorphic forms which may differ from each other in terms of stability, physical properties, spectral data and methods of preparation. However, there is still a need for novel crystalline forms, which are more stable, reproducible and free of other polymorphic forms.\nIn view of the above the present inventors have now found a polymorph of Brexpirpazole hydrochloride, which is stable, having high purity and can be used as an intermediate in the preparation of Brexpiprazole.\nFurther the present inventors have observed that by incorporating these prior art processes the obtained yields of intermediates are low with the formation of by-products which thereby led to the low yields and low purity of Brexpiprazole.\nFurther there is a need for the improved, cost-effective, industrially applicable process for the preparation of Brexpiprazole having high yields and purity.\nThe present inventors have now found a process for preparation of Brexpiprazole, which resolves above referred issues and industrially suitable."}
{"text": "Single wall carbon nanotubes (SWNTs) may be metallic or semiconducting, and may have varying diameters and lengths. Methods of forming SWNTs may produce mixtures of metallic and semiconducting SWNTs with a range of diameters and lengths. Applications using SWNTs may be improved if SWNTs of uniform conductivity are provided, for example all semiconducting SWNTs or all metallic SWNTs. Accordingly, a method and apparatus for sorting SWNTs is desired."}
{"text": "1). Field of the Invention\nThis invention relates to an apparatus that is used for full-wafer testing and/or burn-in testing and/or built-in self-testing.\n2). Discussion of Related Art\nMicroelectronic circuits are usually fabricated in and on semiconductor wafers. Such a wafer is subsequently “singulated” or “diced” into individual dies. Such a die is typically mounted to a supporting substrate for purposes of providing rigidity thereto and electronic communication with an integrated or microelectronic circuit of the die. Final packaging may include encapsulation of the die and the resulting package can then be shipped to a customer.\nIt is required that the die or the package be tested before being shipped to a customer. Ideally, the die should be tested at an early stage for the purposes of identifying the defects that occur during early stage manufacturing.\nThe earliest stage that a die can be tested is after completion of the manufacture of circuits at wafer level and before a wafer is singulated. Full-wafer testing carries with it a number of challenges. One challenge in full-wafer testing is that there are a large amount of contacts on a wafer and that a large number of power, ground, and signal connections thus have to be made."}
{"text": "A variable capacitance diode termed as a varicap diode or a varactor diode is often used to vary a capacitance by controlling a voltage applied to the diode. When a reverse bias voltage is applied, the variable capacitance diode is able to vary the junction capacitance by varying the thickness of the depletion layer in the p-n junction. The variable capacitance diode is highly reliable because the variable capacitance diode includes no mechanical elements. For this reason, variable capacitance diodes are widely used in voltage-controlled oscillators and voltage-controlled variable filters. Variable capacitance diodes are components essential for television receiver sets, cellular phones and the like.\nA variable capacitance diode, however, has a low Q factor (quality factor), and the Q factor largely varies depending on a reverse bias voltage applied to the variable capacitance diode. The Q factor is decreased depending on internal resistance of the variable capacitance diode.\nIn addition, a variable frequency resonator generally uses a variable capacitance diode as a capacitor of the resonator which includes an inductor and the capacitor. Decrease in the Q factor of the variable capacitance diode results in decrease in the Q factor of the resonator. When a resonator is used in an oscillator, the resonator with a lower Q factor deteriorates the phase noise characteristic of an oscillated signal. For this reason, resonators in use for measurement and communications, for instance, are required to have a higher Q factor.\nTwo open-end strip line resonators having different electrical lengths are known to be connected in parallel to a bias supply point of a variable capacitance diode in order to achieve satisfactory phase noise characteristic, for example (see Japanese Patent Application Publication No. 2008-277938)."}
{"text": "In a datacenter, network modules (e.g., Ethernet switches, Virtual Connect fabrics, etc.) may be linked to enhance connectivity and redundancy. Network modules may be linked in a variety of topologies. A ring topology provides a redundant path between each network module connected in a ring. Accordingly, a network module may be connected between two other network modules using two separate communication links. In some examples, the two communication links may send data, messages, packets, etc. in opposite directions (e.g., clockwise and counterclockwise) around the ring.\nWherever possible, the same reference numbers will be used throughout the drawing(s) and accompanying written description to refer to the same or like parts."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing a GaN-based film capable of obtaining a GaN-based film having a large main surface area and less warpage.\n2. Description of the Background Art\nA GaN-based film is suitably used as a substrate and a semiconductor layer in a semiconductor device such as a light emitting device and an electronic device. A GaN substrate is best as a substrate for manufacturing such a GaN-based film, from a point of view of match or substantial match in lattice constant and coefficient of thermal expansion between the substrate and the GaN-based film. A GaN substrate, however, is very expensive, and it is difficult to obtain such a GaN substrate having a large diameter that a diameter of a main surface exceeds 2 inches.\nTherefore, a sapphire substrate is generally used as a substrate for forming a GaN-based film. A sapphire substrate and a GaN crystal are significantly different from each other in lattice constant and coefficient of thermal expansion.\nTherefore, in order to mitigate unmatch in lattice constant between a sapphire substrate and a GaN crystal and to grow a GaN crystal having good crystallinity, for example, Japanese Patent Laying-Open No. 04-297023 discloses forming a GaN buffer layer on a sapphire substrate and growing a GaN crystal layer on the GaN buffer layer, in growing GaN crystal on the sapphire substrate.\nIn addition, in order to obtain a GaN film less in warpage by employing a substrate having a coefficient of thermal expansion close to that of GaN crystal, for example, Japanese National Patent Publication No. 2007-523472 (corresponding to WO2005/076345) discloses a composite support substrate having one or more pairs of layers having substantially the same coefficient of thermal expansion with a central layer lying therebetween and having an overall coefficient of thermal expansion substantially the same as a coefficient of thermal expansion of GaN crystal."}
{"text": "The invention relates generally to folding handle operated door locks, and more specifically to folding T-handle operated door locks easily assembled from a relatively small number of components, which require minimum welding, and which positively open and close.\nLocking door latches are well-known and are used in a variety of applications, including applications where flush mounting is desirable, such as door locks for equipment cabinets and for delivery vehicles. One feature common to locking folding T-handle door latches is that the door latches have a folding operating handle which can be folded down into a recess when not in use. Some drawbacks of prior art door latches include numerous moving parts and/or the requirement of utilizing a number of welds, springs and small parts which can be damaged or rendered inoperable by abuse and/or during painting, and accordingly may be subject to reliability problems.\nThe present invention overcomes drawbacks of prior folding handle latches by providing an improved, relatively simple and inexpensive to manufacture folding handle operated door lock which has relatively few parts, which requires few welds, has no springs or other small parts which can be jammed with painting, and which utilizes a key locking system which could easily retain rotatable components in pre-determined positions relative to housing, namely, a latched and unlatched position.\nThe operating T-handle is preferably nestable within a pan when the T-handle is in its locked position. A key lock is included to releasably lock the T-handle and a latch cam connected to the T-handle. The locking cam and the key lock are spaced apart by virtue of a slide bar which is movable between a locked position and an unlocked position, the locked position being when it is moved into contact with a cutout area of a locking cam connected together with the latch cam. In the unlocked position, the slide bar is moved out of the cutout area of the locking cam. No springs are required to move the slide bar between the locked and unlocked positions, as the action is positive."}
{"text": "Audience measurement of media content (e.g., broadcast television and/or radio, stored audio and/or video content played back from a memory such as a digital video recorder or a digital video disc, audio and/or video content played via the Internet, video games, etc.) often involves collection of content identifying data (e.g., signature(s), fingerprint(s), embedded code(s), channel information, time of consumption information, etc.) and people data (e.g., identifiers, demographic data associated with audience members, etc.). The content identifying data and the people data can be combined to generate, for example, media exposure data indicative of amount(s) and/or type(s) of people that were exposed to specific piece(s) of media content.\nIn some audience measurement systems, the collected people data includes an amount of people being exposed to media content. To calculate the amount of people being exposed to the media content, some measurement systems capture a series of images of a media exposure environment (e.g., a television room, a family room, a living room, etc.) and analyze the images to determine how many people appear in the images at a particular date and time. The calculated amount of people in the media exposure environment can be correlated with media content being presented at the particular date and time to provide exposure data (e.g., ratings data) for that media content."}
{"text": "1. Field of the Invention:\nThis invention relates to a locking mechanism secured to the bottom end of a tubing string which is to be secured in a hollow bore sump packer located at a position in a subterranean well, such as a substantially horizontal portion of the well which traverses a production formation.\n2. Summary of the Prior Art:\nIn recent years, considerable interest has been expressed in utilizing the economies associated with drilling a subterranean well along a path which traverses the length of the formation, rather than the height or depth of the formation. Such drilling necessarily requires that the initial or entry portion of the subterranean well be substantially vertical and then be curved to produce a deviated portion which traverses the production formation and may lie in a substantially horizontal plane.\nThe insertion of a tool string through the curved portion of such deviated well requires the application of a substantial downward force accompanied by a rotation of the tubing string and necessarily generates very high frictional resistance to both movements. As a result, when the end of the tubing string reaches a conventional hollow bore sump packer and is inserted therein, an upward pull on the order of 100,000 lbs. must be exerted on the tubing string to insure that the tubing string is anchored to the sump packer.\nConventional locking mechanisms for sump packers normally involves latching elements which are releasable from the sump packer by the exertion of upward forces thereon on the order of 10 or 15 thousand lbs. The frictional forces acting on the tubing string by virtue of passage through the curved portion of the well bore are typically of a magnitude to make detection at the surface of the well of such a small force change extremely difficult if not impossible.\nThere is, therefore, a definitive need for a latching mechanism for effecting the locking of a tubing string to a sump packer located in a well bore which will withstand, when locked, an upward tensile force on the order of 100,000 lbs., yet may be released from the sump packer through further manipulation of the tubing string.\nAn additional problem arises in the utilization of conventional latching mechanisms in that the position of the locking elements is generally controlled by a J-slot and pin connection between an outer sleeve element of the latching mechanism and an inner mandrel which is secured to the tubing string. Because rotation of the tubing string is required to effect the insertion of the tubing string through the curved portion of a deviated well bore, conventional shear screws cannot be employed to effect the securement of the J-sleeve and locking pin in a desired position during insertion in the well. There is, therefore, a further problem that has to be solved in providing a J-pin and slot mechanism for effecting the movement of the locking elements to and from a locking position which will not be secured by shear screws or the like during insertion of the locking mechanism through the curved portion of the deviated well so that the exact position of the latching elements relative to the sump packer when such elements reach the sump packer is unknown, yet the reliable functioning of the latching elements solely by manipulation of the tubing string is an essential requirement.\nThe principal object of this invention is to provide a locking mechanism for a sump packer located in the deviated portion of a subterranean well which will overcome the above mentioned deficiencies of prior art locking mechanisms."}
{"text": "The present invention relates in general to the field of patient supports, such as hospital beds. In particular, the invention relates to an improved control system that provides enhanced control of the patient support and accessories mounted at the patient support and, further, provides enhanced information related to the patient support and the patient and, further, allows for improved monitoring of the patient."}
{"text": "1. Field of the Invention\nThe invention relates to an electronically commutated DC motor for a liquid pump, comprising a pump housing having a suction connection and a pressure connection, a substantially disk-shaped pump rotor mounted rotatably in the pump housing, the rotor comprising an impeller having a plurality of pump vanes and a permanent magnet, a partition separating a pump space from a dry space, wherein the partition is disposed in an axial gap between the pump rotor and a plurality of axially parallel oriented stator poles of the DC motor, which are each provided with an insulating body and a stator winding, the stator poles adjoin a stator yoke, and the stator containing the stator poles, the stator windings and the stator yoke is disposed in a motor housing made of plastic material.\n2. Description of Related Art Including Information Disclosed Under 37 CFR §§1.97 and 1.98\nAn electronically commutated DC motor of the type in question is known from US 2010/0158724 incorporated by reference herein. It comprises a plastic housing, so that electromagnetic radiation is emitted virtually unimpaired by the stator windings.\nIt is known to use electric motors comprising metal housings to shield against interference. Metal housings, however, are disadvantageous with respect to noise emissions, notably structure-borne noise transmission."}
{"text": "1. Field of the Invention\nThe present invention relates to a toner used in an electrophotographic method, an image-forming method for actualizing an electrostatic image, and a toner jet.\n2. Description of the Related Art\nImage-forming apparatuses using an electrophotographic method have recently been required to demonstrate faster speeds and higher reliability.\nIn addition, there also growing demands for improved appearance of the image as well as energy savings, shortening of wait time and improvement of image productivity, and toners are required to demonstrate superior low-temperature fixing performance in order to respond thereto.\nIn general, low-temperature fixing performance is related to the viscosity of the toner, and toners are required to have the property of being rapidly melted by heat during fixation (so-called sharp melt property).\nAccompanying recent increases in printer speed in particular, the amount of time during which the toner and paper or other media pass through the nip of the fixing apparatus is becoming shorter each year. In addition, there are a growing number of opportunities for users to continuously output images of high quality having a high coverage rate onto heavy paper in the manner of graphic images or posters using an image-forming apparatus such as a laser printer.\nIn the case of using heavy paper, since the amount of heat lost to the paper during fixation increases, it becomes difficult to transfer heat from the fixing member on the paper side to the toner, thereby resulting in inadequate melting of the toner contacting the outermost layer of the paper and increased susceptibility to poor toner fixing performance. In addition, in the case of images having a high coverage rate, since a state results in which much of the toner is layered on the paper during fixation, it again becomes difficult to transfer heat from the fixing member on the toner side to the outermost layer of the paper, resulting in even an even greater likelihood of poor fixing performance.\nConsequently, in the case of outputting images having a high coverage rate using heavy paper, there were many cases in which means were devised for changing the fixing conditions, such as by lowering the printing speed of the fixing apparatus to a speed lower than that of ordinary paper or by increasing the set temperature of the fixing apparatus in order to increase the amount of heat transferred to the toner.\nHowever, image productivity is sacrificed in the case of means for lowering the printing speed of the fixing apparatus, while it becomes difficult to save energy or shorten wait time in the case of means for increasing the set temperature of the fixing apparatus.\nAs has been described above, although there is a need for a toner that can be fixed under the same fixing conditions as ordinary paper even when using heavy paper, a toner capable of satisfying that requirement has yet to be obtained.\nNumerous toners have been proposed as examples of technology for improving toner low-temperature fixability that contain not only an amorphous resin but also a crystalline polyester resin for use as a binder resin (see, for example, Japanese Patent Application Laid-open No. 2003-337443).\nToners containing a crystalline polyester resin are able to improve low-temperature fixability as a result of the crystalline polyester resin melting and becoming compatible with amorphous resin due to the heat during fixation, and the binder resin being plasticized due to this compatibility, thereby resulting in an enhanced sharp melt property.\nHowever, in order to print images having a high coverage rate using heavy paper at high speed, the speed at which the crystalline polyester resin plasticizes the binder resin during fixation (to also simply be referred to as the plasticizing speed during fixation) is inadequate, thereby resulting in the need for further improvement.\nIn addition, when low-temperature fixability is attempted to be improved, there are cases in which toner durability and the rate at which charging rises up become poor, and a toner that satisfies all of these requirements with respect to low-temperature fixability, durability and the rate at which charging rises up has yet to be obtained at the present time.\nFor example, if compatibility between a crystalline polyester resin and an amorphous resin is made to be excessively high in order to improve low-temperature fixability, the crystalline polyester resin ends up melting even at normal temperatures, resulting in a toner that contains plasticized, soft toner particles.\nAs a result, the toner has weak durability with respect to external stress such as that applied when stirring the developer, and in the case of outputting low coverage images such as half-tone images in a mode that is severe on toner deterioration in the manner of continuous output, there is increased susceptibility to a decrease in image density caused by increased attachment force of the toner surface caused by embedding external additives. In this manner, it was difficult to realize both toner low-temperature fixability and durability in toners containing a crystalline polyester resin.\nIn contrast, it has also been considered to enhance toner durability by providing an annealing step in the toner production process and inhibiting compatibility of a crystalline polyester resin by promoting crystallization of the crystalline polyester resin (see, for example, Japanese Patent Application Laid-open No. 2010-152102).\nAlthough it is true that crystallization of a crystalline polyester resin can be promoted by an annealing step, due to the slow nucleation rate during crystallization, crystallization proceeds while the crystalline polyester resin aggregates, thereby resulting in an increased likelihood of the crystalline polyester resin being in a poorly dispersed state.\nDue to the effects thereof, the charge on the surface of toner particles becomes disproportionate during triboelectric charging of the surface, the rate at which charging of the toner rises up ends up decreasing, and there may be increased likelihood of image fogging particularly in the case of continuous output of images having a high coverage rate.\nIn other words, since it was difficult for toners containing a crystalline polyester resin to realize both crystallinity and dispersibility of the crystalline polyester resin, it was also difficult to realize both toner durability and the rate at which charging rises up.\nFurthermore, proposals have also been made to increase the crystallization rate of crystalline polyester resins by adding an inorganic crystal nucleating agent such as silica (see, for example, Japanese Patent Application Laid-open No. 2007-033773) or an organic crystal nucleating agent such as fatty acid amide (see, for example, Japanese Patent Application Laid-open No. 2006-113473).\nHowever, even if these crystal nucleating agents are added, since opportunities for contact with the crystal nucleating agent are limited, crystalline polyester resin ends up remaining that has not crystallized as a result of not being acted on by the crystal nucleating agent, thereby limiting the effect of improving toner durability. In addition, the dispersibility of the crystalline polyester resin was also unable to be improved and the rate at which charging rises up easily became worse. On the other hand, proposals have also been made for improving the amorphous resin combined with the crystalline polyester resin. For example, the use of a hybrid resin containing amorphous resin in the form of a polyester unit and vinyl copolymer unit was proposed in order to improve low-temperature fixability (see, for example, Japanese Patent Application Laid-open No. 2003-173047).\nHowever, due to the inadequate plasticization speed during fixation, further improvements were required for high-speed printing of images having a high coverage rate when using heavy paper.\nIn addition, since improvement of the crystallinity and dispersibility of the crystalline polyester resin in toner was inadequate, there were also problems with toner durability and the rate at which charging rises up.\nAs has been described above, a toner has yet to be obtained that is able to satisfy all requirements relating to low-temperature fixability, durability and the rate at which charging rises up during use of heavy paper."}
{"text": "1. Field of the Invention\nThis invention relates to a nonvolatile semiconductor memory configured by a plurality of nonvolatile memory cells with the transistor structure having a charge storage layer and a manufacturing method thereof.\n2. Description of the Related Art\nRecently, much attention has been paid to a NAND flash memory configured by a series connection a plurality of nonvolatile memory cells with the transistor structure as one nonvolatile semiconductor memory device configuration. It is required for the memory cells (cell transistors) used in the NAND flash memory to satisfy the following two conditions. One condition is to suppress the short channel effect and maintain the preferable turn-on/off characteristic of the transistor even when the gate length is made smaller. The other condition is to raise the channel potential of a non-selected cell to prevent erroneous programming at the programming time.\nIn order to suppress the short channel effect, it is necessary to increase the dopant impurity concentration of the silicon substrate and reduce the depletion layer width. On the other hand, a rise in channel potential required for preventing erroneous programming is realized when the impurity concentration of the substrate is low and the depletion layer width is large. Therefore, as the gate length is reduced, the contradiction between the conditions of suppression of the short channel effect and prevention of erroneous programming becomes more significant. That is, it becomes extremely difficult to miniaturize memory cells while maintaining the basic function of the memory cell.\nAs one solution for the above problem, formation of memory cells on a silicon-on-insulator (SOI) substrate is provided. However, since the procurement cost for the SOI substrate is high, there occurs a problem that the SOI substrate is not suitable for a flash memory in which greater importance is attached to the manufacturing cost.\nThe structure in which thin insulating films are formed under source/drain regions is proposed in order to suppress the short channel effect of a MOS transistor (JP-A H5-218417 [KOKAI]). However, this is the normal transistor structure and is not the transistor structure that realizes a memory cell. Further, even if the structure is applied to a memory cell, the effect of preventing erroneous programming cannot be attained. This is because the insulating films under the source/drain regions are thin and the effect of extending the depletion layer in a depth direction of the substrate cannot be attained.\nFurther, in a NAND cell unit, the structure in which a thin insulating film is formed under a channel region is proposed (JP-A 2007-329366 [KOKAI]). However, the structure is made to form a NAND cell unit with a partial SOI structure and the sufficiently large effect of suppression of the short channel effect and prevention of erroneous programming is not attained."}
{"text": "Metal-insulator-metal (MIM) capacitors are often used in very-large-scale integration (VLSI) designs implemented on semiconductor chips, either as part of the primary circuit design or to provide decoupling capacitance for supply noise suppression. MIM capacitors are intended to improve capacitance density of the VLSI design while minimizing area overhead. FIG. 1 is a cross-section of an example semiconductor chip 100 that includes a MIM capacitor. The MIM capacitor comprises a capacitor top metal (CTM) node 104 and a capacitor bottom metal (CBM) node 106 formed on the top and bottom side, respectively, of a dielectric layer 108. The MIM capacitor layers reside between two metal layers MX and MX+1, on which die interconnection paths 110 are patterned. The MIM capacitor is insulated by an insulating layer 102. Conductive vias 112 are formed in the die to electrically connect circuit paths in the metal layers to the CTM node 104 or the CBM node 106 as needed.\nThe above-described description is merely intended to provide a contextual overview of current techniques and is not intended to be exhaustive."}
{"text": "Food product slicers generally include a food product carriage which is mounted on a slide rod within the slicer housing to allow the carriage to be moved back and forth past a rotating slicing knife. Manual, automatic and combination slicers are known. In the case of automatic or combination slicers a variety of drive arrangements for the food product carriage are also known. In some slicers a multi-link drive arrangement is provided between a rotating motor output and the food product carriage, where the motor rotates in a single direction and change in direction of the food product carriage is achieved mechanically via interaction of the links making up the multi-link drive arrangement.\nIt would be desirable to provide a simple, inexpensive system and method for providing increased control capabilities in connection with such multi-link drive arrangements."}
{"text": "This invention relates to a rubber elastomer composite material having chemical inertness, non-stickiness and low frictional resistance and a process for the production thereof.\nRubber elastomers have been heretofore widely used in a variety of fields of the production of sealing materials, belts, valves, pumps and so on, but have a disadvantage in that they tend to become sticky when used for a long period of time, especially under pressure. For example, it is often observed that a bottle corked for a long period of time with a bottle stopper made of a rubber elastomer is hard to be uncorked, and when the stopper is forcibly pulled out, a part of the rubber is likely to stick and remain on the inner faces of the neck of the bottle. In another example, when a rubber elastomer applied as a valve member or diaphragm for a pump or valve is pressed against metal faces for a relatively long period of time, the valve member or diaphragm is hard to be separated immediately from the metal faces when the pressure is released. This results in a considerable delay in response rate of the pump or valve. Rubber elastomers have other drawbacks in that they have not sufficient resistances to chemicals and abrasion."}
{"text": "1. Field of the Invention\nThe present invention relates to an intake system for an internal combustion engine having a cantilevered valve for opening and closing an intake passage that communicates with an intake port of the engine.\n2. Description of Related Art\nConventionally, an intake system of an internal combustion engine that curbs dispersion of combustion by stabilizing generation of a swirling flow (tumble flow) in a combustion chamber of the engine by reducing a passage sectional area of an intake passage, through which intake air is supplied to the combustion chamber, is proposed (see, e.g., JP-A-2008-25363). As shown in FIG. 10A and FIG. 10B, this intake system of the engine includes an intake manifold 101 of the engine, an intake air flow control valve for opening and closing an intake passage 102 formed in the intake manifold 101, and a rotatable shaft 104 that supports a valve 103 which is a valve body of the intake air flow control valve.\nA cantilevered valve which offsets a centre position of the rotatable shaft 104 (rotation center position of the valve 103) with respect to a centre position of the valve 103 is used as the valve body of the intake air flow control valve. An actuator that includes, for example, a motor for driving the rotatable shaft 104 of the valve 103 and a return spring for urging the valve 103 to be in a fully open state (fully open position) is connected to the intake air flow control valve. The motor is configured such that opening and closing timing of the valve 103, for example, is controlled by an engine control unit (ECU). To fully close the valve 103 of the intake air flow control valve, the rotatable shaft 104 is driven using driving torque of the motor through the supply of electric power to the motor so that a degree of opening of the valve 103 is in a fully closed state (fully closed position). To fully open the valve 103 of the intake air flow control valve, by stopping the supply of electric power to the motor, the return spring urges the valve 103 using spring force of the return spring such that the degree of opening of the valve 103 is in the fully open state (fully open position).\nHowever, the cantilevered valve which offsets the centre position of the rotatable shaft 104 with respect to the centre position of the valve 103, is employed for the intake air flow control valve described in JP-A-2008-25363. In this case, due to differential pressure between upstream and downstream sides of the valve 103 generated when the valve 103 of the intake air flow control valve is fully closed, rotary torque (bending moment) in a direction of valve opening operation with the rotatable shaft 104 of the intake air flow control valve as its center, is applied to the valve 103. Accordingly, maintenance accuracy of the opening degree in maintaining the valve 103 of the intake air flow control valve in the fully closed state (fully closed position) deteriorates. As a result, the motor requires excessive holding torque to maintain the valve 103 of the intake air flow control valve in the fully closed state (fully closed position). Therefore, defects such as upsizing of the actuator including the motor and increase in power consumption of the motor are caused.\nIn addition, the valve 103 of the intake air flow control valve is disposed normally in the intake passage 102 to whose interior portion large negative pressure of intake air and small atmospheric pressure are repeatedly applied in accordance with ascent and descent of a piston of the engine and opening and closing of an inlet valve. Thus, pressure of intake air is applied in a pulsatile fashion to the valve 103 disposed in the intake passage 102. When the pressure of intake air is applied in a pulsatile fashion in the above-described manner, the valve 103 flaps in the intake passage 102. Accordingly, it is difficult to maintain the valve 103 of the intake air flow control valve in the fully closed state (fully closed position) or to maintain the valve 103 of the intake air flow control valve in the fully open state (fully open position), and the opening degree maintenance accuracy in maintaining the valve 103 of the intake air flow control valve in the fully open position and in the fully closed position deteriorates. As a result, the generation of a swirling flow (tumble flow) in a combustion chamber of the engine becomes unstable. Consequently, an effect of improving combustion efficiency of the engine and a fuel efficiency improvement effect by the stabilization of combustion are reduced."}
{"text": "For sampling or dispensing respectively of fluid samples, such as in particular of fluids in the analytical field such as in particular medical or pharmaceutical field among others so called pipettes or capillaries or multiple pipette microstructures are used. In the course of increasing the efficiency in analytical laboratories, the economization and due to smaller sample amounts in this field also the used instruments such as in particular pipettes or multiple pipette structures are getting finer and more complex. For this purpose e.g. from the firms Zymark or Caliper multiple pipette structures are offered comprising up to 384 so called pipette tips for sampling amounts in the range of 2 to 100 nl.\nIn the EP 1 388 369 micro fluidic systems are proposed, which can be used in micro array systems including channel spring probes, that include at least one capillary channel. The proposed spring beams with incorporated channels as e.g. stressy metal beams curve away from a substrate when released. The channel spring probes are arranged onto a substrate by using specific production steps by covering the substrate with a plurality of layer coatings, so that the spring properties are achieved.\nThe subject of the present invention consists in proposing further refinement or increase of capacity and the possibility of increased automation in micro-technical process steps by using Nano- or Microsystems as in particular in analyses, the execution of test series, sampling, sample dispensing, at capillary-electro phoresis, capillary-chromatography, etc., etc.\nAccording the present invention it is proposed, that instead of a “straight” device, a planar 2D-structure or channel spring probes arranged onto a substrate as proposed within the EP 1 388 369, in the nl or μl-range, a bent or arcuate fluidic device or a corresponding fluidic instrument respectively or a 3D-structure is used. E.g. an out of plane sampling device for sampling, transporting and/or disposal of fluid media is proposed, comprising within a substrate plane at least one longitudinal extending strip like portion, comprising a liquid channel, such as a pipette or a capillary tube or needle respectively for sampling or dispensing, which is bent or arcuate designed at least at one location extending out of the plane of the substrate.\n“Bent” in the sense of the present invention means that out of an initially, essentially plane substrate, e.g. a substrate pre-structured by etching a three-dimensional structure has been produced by a specific bending action, the bent parts of the initially plane substrate project out of the plane of the substrate. In particular capillary channels or grooves can be produced which either are running into the substrate plane or also out of the substrate plane.\nSurprisingly and against any assumption it could be proved in experiments that the flow in e.g. so called capillary grooves and in particular in open fluid channels also works “around the corner” in the μl and nano litre range. The open channel or the open channels can be inside as well as 10 outside of the radius of curvature at the location of the bending.\nThe same of course is valid for structures, which consist of a plurality of devices or instruments respectively including liquid channels in the micro litre or nano litre range comprising at least one bent location or structure respectively as proposed above.\nIt goes without saying that the flow characteristic of the fluid within the capillary grooves or fluid channels is dependent on the geometry and the surface finish or coating of the inner wall surface. In case of a water based solution or a fluid the surface is preferably hydrophilic and case of a more oily fluid the surface characteristic is preferably more hydrophobic like. As a general comment, one can say, that the contact angle between the fluid and the surface should be small.\nWithin the prior art so called 3D-structures are known, where capillaries for sampling are bonded into and/or arranged onto a so called “out-of-plane” structure as e.g. the channel spring probes as proposed in the EP 1 388 369.\nFor these structures the production costs are very high and also the production method is complicated and error-prone. In addition with some of such structures closed liquid channels have to be used which means within the prior art often so called closed capillaries are proposed. In addition for the spring probes as proposed within the EP 1 388 369 multi structures have to be applied onto a substrate to achieve the spring properties.\nIn contrast to the today usually used materials for the production of pipettes or capillaries i.e. for the production of the devices proposed according the present invention or structures comprising a plurality of devices or instruments respectively preferably plastically formable material as metals, at least partially plastic polymers and the same are used usually comprising only one layer.\nOn one side the production of the inventively proposed elements or instruments such as pipettes and capillaries or of whole structures is very simple, as the elements or instruments comprising the fluid channels usually consist only of one layer and can be bent in a simple manner. In addition a metal band can be used as a basis, which can be treated by using well-known lithography processes such as e.g. by means of etching to produce the liquid channels. It is possible e.g. to produce open channels on one side or both sides of a small metal band and further more closed channels can be produced by covering the open channels using a film. The possible production processes for the production of the inventive elements or devices respectively like i.e. etching, punching, bending, etc. shall be explained later on with reference to the attached drawings.\nThe advantage of the instruments or devices proposed according the present invention is, that sampling or dispensing of a sample is far easier by using a bent or arcuated element, such as a pipette or a capillary including instrument, as it does not have to be moved against the substrate surface in a more or less perpendicularly manner, from which substrate a sample has to be removed. It is e.g. possible to move the pipette or capillary including device from one side more or less in a parallel manner over the surface of the substrate for removing a sample. By using an inventive structure also at limited dimensional conditions at the same time a plurality of samples can be removed or dispensed in a simple manner. The various advantages of the inventively proposed devices or elements respectively as well as the structures shall be explained in more details later on with reference to the attached drawings.\nFurthermore a process is described for sampling, transporting and/or dispensing of fluid media in particular by using the above mentioned devices or the above mentioned structures. According to the inventive process a pipette or capillary such as a needle or a structure, comprising a plurality of pipettes or capillaries is used and the fluid is transported along a location at sampling or dispensing a sample, which is arcuated or bent respectively.\nFurther preferred aspects or embodiments of the elements, devices as well as of the structures and the processes respectively are characterized in dependent claims.\nThe devices, instruments or structures respectively according the present invention are suitable in particular for diagnostic or analytical processes in the field of chemistry, medicine, microbiology, pharmaceutics, etc. The invention shall be explained in more details with reference to the attached drawings."}
{"text": "The invention relates generally to backpacks and, more particularly, to backpacks having a waistbelt.\nBackpacks have long been used for carrying heavy, bulky loads. Over the years, backpacks of various configurations have been made, including packs having external frames, internal frames, and those without frames. External-frame backpacks typically include interconnected metal bars, forming a relatively rigid structure. A pack body, typically of nylon or canvas, is secured generally within the confines of the frame and, as such, is relatively spaced apart from the back of the user. Internal-frame backpacks typically include internal stays disposed within pockets of the body of the backpack, allowing the backpack to be positioned more closely to the back of the user. Frameless backpacks typically exclude rigid support structures, allowing the pack to conform to the user.\nRegardless of type of backpack, much effort has been made to distribute the weight of the load evenly onto the user. The ability to carry heavy loads in relative comfort depends in great part upon the placement and transfer of the load. To that end, backpacks often include a waistbelt to transfer some of the load onto hips and lumbar area of the user. Typically, the belt is attached to a lower portion of the backpack, extending about the hips of the user.\nTraditionally, waistbelts have been provided as a padded belt having an attachment point to the sack of the backpack in a rear portion and a buckle in a front portion. Such belts typically were formed of padding material and fabric and were cinched around the user. The belts aided in distributing load onto the hips. However, as a user moved, the weight tended to shift excessively, overloading concentrated areas on the user, causing discomfort. Moreover, such belts tended to be shaped for a generic body type.\nMore recently, certain waistbelts have been configured with adjustments to better conform to a range of body types to include different hip-to-waist angles. In one example, a belt includes two lateral arms that extend from a support of the backpack and wraps around the user's hips. The arms can be set at three different angle settings to accommodate different hip-to-waist angles, by attaching screws in one of three different positions provided cooperatively by the support and the arms. Once set, the lateral arms are fixed in at that angle. Thus, as the user moves the load still tends to shift relative to the hips, as previously discussed.\nIt should, therefore, be appreciated that there remains a need for a backpack that accommodates wide range of body types and provides effective load distribution. The present invention fulfills this and other needs."}
{"text": "A resonant accelerometer is a sensor that responds to an acceleration force by producing a frequency shifted output signal. Quartz-based resonant accelerometers have been used in many commercial applications, including navigation-grade precision accelerometers.\nMicromachined resonant sensors have been developed. The acceleration force amplification provided by these early devices has been limited by the leverage systems for the proof masses of the devices. Thus, to improve the response of micromachined resonant sensors, it is important to improve upon prior art proof mass leverage systems.\nSome recent work has focused on micromachined resonant sensors in bulk silicon processes, but this class of sensor has not yet been pursued in a surface-micromachining technology. Surface-micromachining technology embeds a micromechanical device in an anisotropically etched trench below the surface of a wafer. Prior to microelectronic device fabrication, this trench is refilled with oxide, chemical-mechanically polished, and sealed with a nitride cap in order to embed the micromechanical devices below the surface of the planarized wafer. The wafer is then used as the starting material for integrated circuit fabrication in a conventional process, such as CMOS or BiCMOS. Thus, surface-micromachining technology allows a micromachined device to be combined with integrated circuitry in a single wafer.\nIn view of the foregoing, it would be highly desirable to provide a resonant accelerometer with an improved leverage system for enhanced force amplification. In addition, it would be highly desirable to provide a resonant accelerometer design that is compatible with surface-micromachining technologies."}
{"text": "The present invention relates to a data retaining circuit such as a latch circuit. More particularly, the present invention relates to a data retaining circuit that reduces the occurrence of soft errors due to alpha rays, neutrons and so forth.\nThe semiconductor devices used in space or in aircraft have a problem of the occurrence of soft errors due to radioactive rays such as alpha rays, neutrons and so forth. Recently, semiconductor devices have been developed into those with higher integration and lower voltages, resulting in a problem that the occurrence of soft errors due to radioactive rays begins to have influence even on semiconductor devices used on the ground.\nFIG. 1 is a diagram that shows the relationship between the amount of critical charge and the soft error rate due to alpha rays and neutrons in a latch circuit that is manufactured in 0.35˜im process and operates at 3.3V. The soft error rate is expressed by 1 FIT=1 error/109 device time. As shown schematically, when the amount of critical charge is equal to or less than 75 fC, the occurrence frequency of soft errors due to alpha rays is high and when the amount of critical charge is equal to or greater than 75 fC, the occurrence frequency of the soft errors due to neutrons is high. The amount of critical charge of a latch circuit that is manufactured in 0.35°m process and operates at 3.3V is approximately 150 fC, therefore, the problem of soft errors due to neutrons is more serious in this circuit. Devices used in space equipment or in aircraft have problems such as an increase in the failure rate of pacemakers in an aircraft. It is obvious that the increase in the occurrence of soft errors due to radioactive rays can be expected also in the semiconductor device used on the ground and the influence cannot be neglected if the structure of a semiconductor device becomes finer in the near future. Moreover, not only neutrons but also alpha rays are expected to bring about a problem of soft errors. With the above-mentioned fact as the background, countermeasures against the soft errors are required in various fields of the semiconductor device.\nAlthough the occurrence of soft errors in a memory can be detected and corrected by providing redundant bits, countermeasures against soft errors are required also in logic circuits. In logic circuits, a data retaining circuit such as a latch circuit has largest influence when soft errors occur therein. Even if the data in such as a combinational circuit is temporarily reversed, it returns to the original one as long as the data in the former stage is not reversed, therefore, the range of the influence can be limited, but if the retained data is reversed and retained as it is the reversed data propagates and produces a far-reaching influence. Because of this, the countermeasures against soft errors in a data retaining circuit are especially required and the present invention relates to countermeasures against soft errors in a data retaining circuit.\nFIG. 2A and FIG. 2B are diagrams that illustrate the mechanism of occurrence of soft errors. As shown in FIG. 2A, a transistor comprises a gate G, a source S and a drain D, and a channel is formed between the source and the drain under the gate and a depletion layer E is formed around the source and the drain. When particles such as alpha rays and neutrons enter the drain region of the transistor, many pairs of holes and electrons are generated on the trajectory of the particles due to the collision with atomic nuclei. At this time, the shape of the depletion layer becomes enlarged by the pairs of hole and electron. This region is called the funneling region F.\nIn the depletion layer and funneling region, electrons and holes move by drifting as shown in FIG. 2B. As movement by drifting is performed at a high speed, one group of generated charges quickly moves in the direction toward the drain, as a result. Electrons and holes generated in other than the depletion layer and the funneling region move by diffusion, but movement by diffusion is performed at a lower speed than that by a drift, therefore, almost all of the pairs of holes and electrons disappear, in pair annihilation, by annihilating each other, but part flows into the depletion layer and the funneling region and moves by drifting in the direction toward the drain.\nThe direction of movement of electrons is different from that of holes in NMOS and in PMOS, that is, electrons flow into the drain node in NMOS and holes do so in PMOS. At this time, when the amount of charge that flows into the drain node is larger than the amount of critical charge of the node, a phenomenon occurs in which the data retained by the transistor is reversed, which is called a soft error. In a MOS transistor, a soft error occurs characteristically in such a way as to reverse from a high logical level (data: 1) to a low logical level (data: 0) in NMOS, and in such a way as to reverse from a low logical level (data: 0) to a high logical level (data: 1) in PMOS because of each structure.\nAs samples of countermeasures against the soft errors in a data retaining circuit, K. Joe Hass, Jody W. Gamples: ‘Mitigating Single Event Upsets From Combinational Logic” 7th NASA Symposium on VLSI Design 1998 has disclosed the circuit as shown in FIG. 3, and U.S. Pat. No. 6,026,011, the circuit as shown in FIG. 4. Each circuit has a structure in which the node that retains the latched data is divided into a portion composed only of NMOS and the other portion composed only of PMOS, each having the same data and correcting the retained data for each other, the characteristic that only errors from 1 to 0 occur in NMOS and only those from 0 to 1 occur in PMOS being taken into consideration. Because each has the same data, soft errors occur only in one of them and the errors do not occur in the other, therefore, the data in which the error occurs is corrected by the data in the other portion where no error occurs.\nTo put it concretely, in the circuit in FIG. 3, when input data D is 0, data PP and NN to be taken in are also 0, and QP and QN become 1 and an output Q becomes 0. In this case, PP retained in the data retaining section composed of a PMOS has a possibility of the occurrence of a soft error of changing from 0 to 1, and QN retained in the data retaining section composed of an NMOS has a possibility of that of changing from 1 to 0, but it is unlikely that NN and QP reverse, When QN changes from 1 to 0, the NMOS transistor of the output section turns off, the PMOS transistor of the output section is off because QP is 1, and the output Q enters a state of floating, but the data does not reverse because of the parasitic capacitance. Then QN is brought back to the original data 1 due to NN and QP, therefore, the normal state returns. As for PP, it is also brought back to the original data 0 due to NN and QP. When the input data D is 1, there is a possibility that a soft error occurs in NN and QP, but since no soft error occurs in PP and QN, the original stage returns similarly.\nIn the circuit in FIG. 4, the data PP and NN are the same as the input data D and, as they are retained in the data retaining section composed of a PMOS and in that composed of an NMOS, respectively, a soft error will occur only in one of them, therefore, it is possible for them to return to the original state because they correct each other. On the other hand, the possibility that data HLD in the output data retaining section changes is very low because the outputs of two inverters, the gate inputs of which are the data PP and NN, are connected commonly and if one of the outputs changes temporarily due to a soft error, the other output is maintained correctly.\nIn the circuits in FIG. 3 and in FIG. 4, when the amount of charge to cause a soft error is small and the reversion of the data due to the soft error is returned to the original state in a brief time by correction, it is possible to maintain a correct state, but if the state of reversion of the data lasts a long time, the reversed data propagate to all over the circuit and the retained data is completely reversed as a result. In this case, the reversed data is retained as is. It is unlikely in actual use that the state of reversion of the data lasts so long that the retained data is completely reversed, and it does not bring any problem from a practical viewpoint. In the future, however, as semiconductor devices become finer, resulting in the reduction in capacitance components and voltages becoming lower, resulting in the reduction in the amount of critical charge, the state of data reversion caused by a soft error becomes longer accordingly and there is a possibility that the occurrence of soft errors cannot be prevented sufficiently even by using the circuits in FIG. 3 and FIG. 4.\nIn the circuit structures in FIG. 3 and FIG. 4, the node of the NMOS side and that of the PMOS side are connected so as to feed back each other, therefore, a problem occurs that the operation speed is slow. Moreover, another problem occurs that the structure is complicated and the size of the circuit becomes larger because many transistors are used.\nThe semiconductor device is required to become more dense in circuit integration, speedier, and more economical in power consumption, and the data retaining circuit such as a latch is in the same situation. It is proposed, therefore, that a pulse latch is used as a flip-flop in a structure as shown in FIG. 5. Such a latch, however, retains data dynamically, therefore, it has a characteristic that it is very weak to a soft error. It is, therefore, required that the data retaining circuit as shown in FIG. 5 has an improved resistance to soft errors."}
{"text": "This invention relates to a headrest of the kind which is movable relative to the backrest of a seat on which it is mounted.\nWith headrests of the above kind, movement usually varies the elevation of the headrest relative to the backrest, but it is also known to have a facility for fore and aft adjustment of the headrest position. It will be convenient to hereinafter describe the invention by reference to an adjustable headrest mounted on a vehicle seat, but it is to be understood that the invention has wider application and can by applied to seats of all kinds.\nIt is known to provide vehicle seats with headrests which are capable of adjustment so that the elevation of the headrest is varied. It is also known to provide headrests with a fore and aft adjustment facility. In some cases headrests are arranged to be adjustable both in relation to elevation and fore and aft disposition. Furthermore, it is known to provide such headrests with drive systems through which adjustment of the headrest position is effected.\nIn U.S. Pat. No. 4,830,434, a drive system is used for adjusting the headrest elevation, but not the fore the aft position. In the arrangements disclosed by U.S. Pat. Nos. 5,052,754 and 5,222,784, a drive system is provided for both types of adjustment. The drive system disclosed by each of the three U.S. patents is relatively complicated.\nVehicle manufacturers typically offer customers a choice of options concerning seating arrangements, including headrest adjustment. Options relating to the headrest may include manual adjustment, power adjustment, two-way adjustment (fore and aft, or up and down), four-way adjustment (both fore and aft, and up and down) and position memory. Given the varying requirements of those options it is generally not possible to offer all options without substantial variation in cost. By way of example, different options may require different seat backrest constructions or different headrest mounting systems.\nThe present invention seeks to provide a headrest which is capable of four-way adjustment and which is amendable to relatively simple construction. At least in preferred forms, the invention seeks to utilise a modular system such that principal components of the headrest can be utilised in a headrest capable of two-way (height only) adjustment.\nAn adjustable headrest according to one aspect of the invention includes a frame structure which can be secured on a mounting system for up and down movement relative to an associated seat backrest, and which includes a height adjuster section and a fore/aft adjuster section. The height adjuster section is secured on the mounting system for relative up and down movement of the frame structure. The fore/aft adjuster section is engaged with the height adjuster section for relative movement providing fore/aft adjustment. The headrest has drive means which includes two drive systems which control height adjustment and fore/aft adjustment respectively. The fore/aft drive system includes an adjustment member which is constrained substantially to up and down movement and which is engaged with the fore/aft adjuster section so as to permit fore/aft movement of the fore/aft adjuster section relative to the height adjuster section and the adjustment nut.\nThus, according to the invention, there is provided an adjustable headrest, for a seat backrest having an associated mounting system, wherein the headrest includes a frame structure adapted to be secured on the mounting system for up and down movement relative to the seat backrest, the frame structure including a height adjuster section and a fore/aft adjuster section; the fore/aft adjuster section is engaged with the height adjuster section for relative movement providing fore/aft adjustment; wherein the headrest further includes drive means which has a height adjustment drive system for controlling height adjustment of the headrest by moving the height adjuster section and fore/aft adjustment drive system for controlling fore/aft adjustment of the headrest by moving the fore/aft adjuster section; and wherein the fore/aft drive system has an adjustment member or nut which is constrained substantially to up and down movement and which is engaged with the fore/aft adjuster section so as to permit fore/aft movement of the fore/aft adjuster section relative to the height adjuster section and the adjustment nut.\nThe fore/aft adjuster section may be engaged with the height adjuster section by first engagement means which precludes relative movement therebetween in a direction for height adjustment, but which enables relative movement therebetween in a direction for fore/aft adjustment. The first engagement means may comprise at least one projection, rib or fin on one of the sections which is located in a slot which is defined by the other section and which extends in the fore/aft direction. The adjustment member or nut may be engaged with the fore/aft adjuster section by engagement means which constrains relative movement therebetween to movement along a direction which is inclined with respect to directions for both height and fore/aft adjustment to provide thereby fore/aft adjustment due to relative movement between the height and fore/aft adjuster sections being precluded in a direction for height adjustment. The second engagement means may be similar to the first engagement, but with the or each slot extending in the inclined direction. In each of the engagement means, the or each slot may be a groove, channel or the like, or a spacing between projections, ribs or fins defined by the component having the or each slot.\nIn one arrangement, the height adjuster of the frame structure is engaged with the fore/aft adjuster by having at least one projecting portion, such as a tab or the like, which engages in or with an elongate slot or track which is defined by the fore/aft adjuster and which extends in the fore/aft direction, i.e. in the direction required for the fore/aft adjustment. Similarly the adjustment nut is engaged with the fore/aft adjuster by having at least one projecting portion, such as a tab or the like, which engages in or with an elongate slot or track defined by the fore/aft adjuster and which extends in the fore/aft direction but is inclined at an angle to the slot or track engaged by the height adjuster. For engagement of one or each of the height adjuster and the adjustment nut with the fore/aft adjuster, a complementary method of this engagement can be used, but this can tend to limit the extent of fore/aft movement or provide a less compact arrangement. Thus, the height adjuster may define at least one track or slot in or with which a projecting portion of the fore/aft adjuster engages, while the adjustment nut may define at least one track or slot in or with which a projecting portion of the fore/aft adjuster engages.\nIn one embodiment, the drive system for fore/aft adjustment is an electrically powered drive means operable to raise and lower the adjustment nut, that is, to move it in the direction for height adjustment, relative to the height adjuster of the frame structure. Preferably the drive means includes an electric motor and a rotatable screw-threaded spindle drivingly connected to the motor, such as through a worm drive. The motor and spindle may be secured on the height adjuster, with the spindle co-operatively engaging a screw-threaded bore defined by or secured to the adjustment nut. The spindle is disposed in the direction of height adjustment such that, with rotation of the spindle by the drive motor, the adjustment nut is caused to travel along the spindle. Because of the respective engagement between the fore/aft adjuster of the frame structure and each of the height adjuster and the adjustment nut, the travel of the adjustment nut causes the fore/aft adjuster to move in the fore/aft direction, relative to the height adjuster, to provide fore/aft adjustment of the headrest.\nThe drive system for height adjustment may be as disclosed in PCT/AU97/00325. Thus, the drive system may include an electrically powered drive means operable to raise and lower the height adjuster of the frame structure, and hence the headrest, relative to the mounting system. Preferably the drive means includes an electric motor and a rotatable screw-threaded spindle drivingly connected to the motor, such as through a worm drive. The motor and spindle may be secured on or in relation to one of the height adjuster and the mounting system, while the spindle may co-operatively engage a screw-threaded nut defined by or secured on or in relation to the other one of the height adjuster and the mounting system.\nThe arrangement for fore/aft adjustment for the headrest of the invention differs from known arrangements, in particular those known arrangements using a tilting or nodding action. The arrangement of the invention enables fore/aft adjustment which is more effective for a greater height range for users. Moreover, throughout that range, the headrest can be maintained in a constant orientation, as the fore/aft adjustment can be and most preferably is substantially linear. For this, the or each groove or track by which the height adjuster of the core structure engages the fore/aft adjuster most preferably is substantially linear, with the groove or track most preferably substantially perpendicular to the direction of height adjustment.\nAs indicated, the or each groove or track by which the adjustment member engages the fore/aft adjuster of the frame structure is inclined with respect to the or each groove or track by which the height adjuster engages the fore/aft adjuster. The inclination may be such that the grooves or tracks converge in the forward direction, such that as the adjustment nut is raised or lowered, the headrest is adjusted respectively in the forward or rearward direction. The angle at which the grooves or tracks converge may, for example, be from about 20xc2x0 to 40xc2x0, and determines the magnitude of fore/aft adjustment per unit length variation in the height of the adjustment nut.\nIn one embodiment the headrest is adapted for use with a mounting system including a post, or each of two laterally spaced posts, secured to the seat backrest and projecting beyond the upper edge of the backrest. For ease of description, the more usual arrangement of two laterally spaced posts is assumed in the following. The headrest frame structure may be mounted on both of the posts for relative sliding movement thereon in the axial direction of the posts. A cross member may extend between and be secured to the upper end portion of each of the posts, in which case the drive system for controlling height adjustment of the headrest may extend between that cross member and the height adjuster of the frame structure.\nA power headrest according to the invention can be provided with a memorised position facility, and/or with a rapid-drop facility. Also, at least one support post of a powered headrest can form or include part of the electrical circuit through which the or each drive system is controlled. Still further, the or each post can form part of two or more individual circuits including (for example), the elevation circuit control, an upper memorised position circuit, and a lower memorised position circuit and/or corresponding fore/aft position circuits."}
{"text": "This invention relates to a semiconductor luminescent device including a control electrode.\nSemiconductor luminescent diodes emit light by a forward current flow through the PN junction involved. In order to control the luminescent output from a luminescent, it is required to externally connect a switch element such as a transistor or a variable impedance element to each semiconductor luminescent diode. Such a switch element has an inherent power loss and is required to have a current rating similar to that of the corresponding semiconductor luminescent diode. Thus the conventional semiconductor luminescent diodes are disadvantageous because the control of their luminescent output is expensive.\nThere have been already known luminescent elements including a control electrode for control of the luminescent output, for example, luminescent switching elements having a four layer PNPN structure and the control or gate electrode disposed in ohmic contact with either one of the intermediate layers to be formed into a thyristor configuration. Once such an element has been turned on with an associated source of direct current connected thereacross, it continues to remain on to emit light unless the current decreases to less than the holding current. This means that the control of the luminescent output requires a transistor or another variable impedance element for externally controlling the on-state or forward current.\nIt is an object of the present invention to provide a novel semiconductor controlled luminescent device.\nIt is another object of the present invention to provide a new and improved semiconductor controlled luminescent device including an extremly simple control circuit.\nIt is still another object of the present invention to provide a new and improved semiconductor controlled luminescent device controlled with low power loss.\nIt is a further object of the present invention to provide a new and improved semiconductor controlled luminescent device having an extremly high control speed."}
{"text": "The invention relates to a selective electronic tripping device associated to a circuit breaker and comprising:\nmeans for measuring the current flowing in the circuit breaker,\nprocessing means connected to the means for measuring the current and comprising means for determining a quantity representative of the peak value of the current, means for comparing said quantity with an instantaneous tripping threshold so as to provide an instantaneous tripping signal when said quantity exceeds said threshold.\nKnown electronic tripping devices comprise a processing unit, generally with a microprocessor, performing in particular long delay, short delay and instantaneous tripping functions. To achieve selectivity between two circuit breakers connected in series or cascading, it is known to provide a high rating for the up-line circuit breaker and a lower rating for the down-line circuit breaker. In this case the instantaneous tripping selectivity is of the current intensity type, the instantaneous tripping threshold of the down-line circuit breaker being lower than that of the up-line circuit breaker.\nIt is also known to use a time-based selectivity between two circuit breakers. In this case the tripping device of the up-line circuit breaker has a greater instantaneous tripping delay than the instantaneous tripping time of the tripping device of the down-line circuit breaker.\nThese two types of selectivity are not always efficient when short-circuit currents are very high, since high current thresholds are reached in both circuit breakers and a too long time delay of the associated tripping device would be liable to lead to premature wear of the contacts of the up-line circuit breaker.\nTo reduce these drawbacks a selective tripping device has been proposed, in particular in the document EP-A-128084, comprising a counter enabling the successive opening and closing cycles of the circuit breaker contacts in case of a short-circuit to be counted.\nIt has also been proposed, in particular in the document EP-A-872939, to disable an instantaneous tripping order if repulsions are detected for primary currents lower than a preset threshold.\nAnother type of selective instantaneous tripping, which will be described in greater detail with respect to FIGS. 3 and 4, uses one time delay when the current has exceeded a first, low, threshold, and a second, higher, threshold after this time delay.\nIn all known selective tripping devices, it is in practice necessary to wait for at least two short-circuit current peaks before making an instantaneous tripping decision. This leads to detrimental wear of the circuit breaker contacts.\nThe object of the invention is to provide a selective tripping device not presenting these shortcomings. The selective tripping device according to the invention has to be faster than known selective tripping devices, while at the same time achieving selectivity with a down-line circuit breaker so as to reduce wear of the contacts as far as possible in the event of a short-circuit.\nThis object is achieved by the fact that the processing means comprise means for determining a second quantity, representative of the time necessary for the current to reach the peak value, and means for determining the instantaneous tripping threshold according to a decreasing function of the second quantity.\nAccording to a development of the invention, the instantaneous tripping threshold varies according to the second quantity according to a curve disposed between first and second curves representative of envelopes of the peak values of the current, the first curve being obtained when the associated circuit breaker is not connected in series with another circuit breaker connected down-line and the second curve being obtained when the associated circuit breaker is connected in series with another circuit breaker connected down-line. The first and second curves are preferably determined experimentally.\nThe tripping device according to the invention thus enables the associated circuit breaker not to be opened for no purpose when the latter is connected in series with a down-line circuit breaker. However, in the absence of a down-line circuit breaker, the absence of the latter is detected in less than one current half-period, which enables the destructive effect of repulsion of the contacts in case of a short-circuit to be minimized."}
{"text": "The present invention relates to medical devices, methods and systems. More specifically, the invention relates to devices, methods and systems for promoting fusion of and/or stabilizing a facet joint of a spine.\nApproximately 80% of Americans experience at least a single episode of significant back pain in their lifetime. For many people, back pain is a chronic, often debilitating disorder. The direct costs of treating back pain, as well as the indirect costs, such as lost wages and decreased productivity, are staggering.\nAlthough back pain may be caused by a number of different factors, many cases of back pain are caused by conditions related to the spinal (or “vertebral”) column. The vertebral column is made up of bones (vertebrae) and intervertebral discs that reside in the joint spaces between the vertebral bodies of the vertebrae. Three joints reside between every two adjacent vertebrae—one larger intervertebral joint between the two intervertebral bodies and two facet joints located posterolaterally relative to the vertebral bodies. These three joints share the load applied between every set of two vertebrae. Many spinal ailments are caused by degeneration, injury or deformity of these vertebral joints and/or intervertebral discs.\nTreatment of spinal pain typically begins with conservative, non-surgical methods, such as rest, heat, analgesics, physical therapy and manipulation. Unfortunately, however, conservative treatments fail in a significant number of spinal pain patients, and surgery is often required. Current surgical procedures for treating intervertebral disc and joint maladies include decompressive surgery, in which all or part of an intervertebral disc and/or the spinal laminae and facets are removed, decompression with fusion of the joint (or “arthrodesis”), and arthrodesis alone. Intervertebral joint arthrodesis involves fusing two adjacent vertebrae to stop the motion between those vertebrae.\nDecompression (removal of structures compressing the spinal nerves including laminae, facet joints and/or intervertebral discs) is a very common surgical procedure that is very effective in promptly relieving significant pain derived from pressure on spinal nerve roots (“radicular” pain). The overall success rates for decompression alone, however, range from 48% to 89%. Spinal fusion is sometimes needed in combination with decompression to more successfully treat spinal pain.\nIntervertebral fusion is designed to stop the motion at a painful vertebral joint, which in turn should decrease pain generated from the joint. It is also performed to stabilize an unstable intervertebral segment, which if left unfused, could cause recurrent compression of the spinal nerves. Fusion procedures involve adding bone graft to an area of the spine to set up a biological response that causes the bone graft to grow between the two vertebral elements and thereby stop the motion at that segment. Often, some type of support structure is attached to the two vertebrae being fused, to hold the vertebrae in a stable position relative to one another while the bone graft material causes fusion. Typical support structures, for example, include bone screws (or “pedicle screws”) attached to rods. Discectomy combined with fusion has been the most common surgical treatment for symptomatic cervical spondylosis for over 40 years. Good to excellent results have been reported in 52-100% of anterior lumbar interbody fusions and 50-95% of posterior lumbar interbody fusions.\nA number of different spinal fusion surgical procedures are currently in use. The most common fusion procedure performed in lumbar surgery is posterolateral intertransverse fusion. Although the technique is often quite successful, the standard procedure often causes significant trauma to the paraspinous muscles. These muscles must be stripped from the transverse processes and retracted for an extended period of time to expose the underlying bone, which can result in denervation, devascularization, and ischemia of these important muscles, leading to atrophy and necrosis. Significant spasm in the short term and atrophy and necrosis in the long term contribute to the morbidity and sequelae of fusion. The clinical effect of this muscle morbidity can be significant postoperative pain and functional impairment in the convalescent period, as well as permanent impairment of paraspinal lumbar muscular function.\nOther surgical techniques for performing spinal fusion are associated with similar and/or additional risk factors. For example, posterior lumbar interbody fusion (“PLIF”) achieves fusion by inserting bone grafts, titanium threaded cages, bone dowels, or carbon fiber spacers filled with bone graft into an intervertebral disc space. All PLIF techniques require removal of the posterior bone of the spinal canal (laminectomy), retraction of the nerves and removal of the disc material from within the disc space, any of which may cause complications. Another procedure, anterior lumbar interbody fusion (“ALIF”), is similar to PLIF, except that in ALIF the disc space is fused by approaching the spine through the abdomen instead of through the back. An additional, and potentially significant, risk of ALIF is potential damage to abdominal structures, such as the large arteries that supply blood to the legs.\nAs mentioned above, with many spinal fusion procedures, some type of support structure, such as screws, rods, pins, cages and/or the like, is used to hold the adjacent vertebrae in place while they are fusing together with the help of the bone graft, bone adhesive, or the like. One of the primary risks of fusion surgery is that a solid fusion will not be obtained (“nonunion”), thus requiring further surgery. One of the main challenges of intervertebral fusion surgery is to stabilize the vertebrae long enough, using the support structure(s), so that they have time to fuse. Another challenge is actually applying the support structure properly. Pedicle screws, for example, may often be effective at providing support during intervertebral fusion, but they can be difficult to place properly, and if misplaced may cause nerve root and/or vascular injury.\nAs the population continues to age, surgical procedures for fusing and/or stabilizing vertebrae will become ever more common. Therefore, a need exists for improved techniques, devices and systems for performing such procedures. Such improved methods and devices should ideally facilitate and/or enhance intervertebral fusion, while preventing or reducing the prevalence of complications or sequelae. Ideally, minimally invasive procedures would be developed that would provide stabilization of an intervertebral joint for a sufficient period of time to allow the vertebrae to fuse. Also ideally, such procedures would be relatively simple to use. At least some of these objectives will be met by the present invention."}
{"text": "This invention relates to garbage collectors that scan and summarize heap memory as part of a garbage collection process. In general, memory reclamation may be carried out by a special purpose garbage collection algorithm that locates and reclaims memory which is unused, but has not been explicitly de-allocated. There are many known garbage collection algorithms, including reference counting, mark-sweep, mark-compaction and generational garbage collection algorithms. These, and other garbage collection techniques, are described in detail in a book entitled “Garbage Collection, Algorithms For Automatic Dynamic Memory Management” by Richard Jones and Raphael Lins, John Wiley & Sons, 1996.\nAn object may be located by a “reference”, or a small amount of information that can be used to access the object data structure. Garbage collectors typically examine these references to determine which objects are reachable and, thus, must be maintained and which objects are no longer reachable and, thus, constitute garbage that can be reclaimed. However, when application threads run they can create, delete and modify references. Accordingly, some mechanism must be used to inform the collector of changes made to references by application threads while the collector is examining the references.\nWrite barriers are one method used to notify garbage collectors of changes in reference values in the heap caused by application threads. Typically, a write barrier intercepts an attempt by an application thread to store a value to memory. The write barrier can then check to determine whether a reference is being modified and, if so, mark the reference as modified, inform the collector of the modification or perform other processing that will later inform the collector that the reference has been modified. Write barriers are discussed in general in the aforementioned book by Jones and Lins.\nWrite barrier techniques have typically fallen into three categories: (1) hardware barriers and memory protection, (2) card tables, and (3) sequential store buffers (SSBs). Hardware barriers and memory protection schemes use special hardware or the memory protection provided by the virtual memory system to detect memory stores by the application threads.\nCard tables are typically used with generational garbage collectors, which have younger generations that are collected frequently and older generations that are collected less frequently. In such systems it is necessary to track references from older generations into younger generations so that the younger generations can be collected without examining every object in the older generations. In a card table technique, the heap memory belonging to the older generations is split into equal-size chunks called “cards.” A card table is an array with a one entry per card in the heap. When a reference update occurs, the card containing the updated reference is marked “dirty” by setting its entry in the card table to an appropriate value. Later, during the collection process, only heap memory areas corresponding to dirty cards are scanned.\nSequential store buffers are thread-local buffers that store the addresses of fields into which cross-generational references are stored. The references may be filtered before storing them to eliminate references that are not of interest and to eliminate duplicates in the buffers. An overflow of a buffer is trapped and causes the buffer to be processed and recycled for further use.\nBecause most collection is driven either by the allocation activity of the application threads or by the collector, if it is performing collections concurrently with the operation of the application threads, one challenge has been how to control how much memory must be scanned and summarized by the collector when it suspends the application threads to perform a collection. The ability to limit the amount of memory that must be scanned is particularly important for space-incremental techniques like the Train algorithm. Some collectors address this problem by limiting the amount of memory the application may mark as dirty, while other collectors have employed concurrent threads to scan modified locations. The card table and sequential store buffer write barrier techniques tend to collapse repeated stores to the same or close reference locations to the same card or buffer. Since only dirty cards or buffers into which field addresses have been stored are scanned, the amount of scanned memory is reduced. However, the use of card tables incurs two costs. First, the number of cards that must be scanned can approach the number of cards comprising the generation with which they are associated; second, all entries of the card table must be examined regardless of which ones are marked as dirty; and, third, when a card is processed, all reference locations in the card must be examined regardless of which ones have been recently modified. As the size of the generation associated with the cards and card table increases both of these costs increase. Sequential store buffers can determine entries to be scanned more exactly, but incur additional processing when a buffer becomes full."}
{"text": "At present, on a common image forming device, such as a printer, a device for accumulating recording medium including papers etc. (in order to facilitate description, recording medium are hereinafter referred to as papers) is commonly installed, which is intended to avoid the printed pages from being taken away by users when the printing has not yet been finished and thus resulting in the loss of a part of the information. A printer with a medium accumulating device will not discharge the printed page from the printer until one sheet of printed page has been finished. However, a conventional printer with a medium accumulating device can only collect one piece of printed page at one time, i.e., after the printing of one piece of printed page has been finished, the piece of printed page will be discharged from the printer. However, when the printed contents are more than one sheet, users are not usually clear about how many pages are left behind, and thus often take the discharged printed pages away with those left behind being omitted, therefore resulting in a disclosure of users' information.\nA device capable of accumulating a plurality of printed papers is disclosed in an international patent published in 2000 (Publication No. U.S. Pat. No. 6,077,030A). Referring to FIG. 1A, the device includes a drum cylinder 20 capable of receiving and accumulating a plurality of papers by means of predetermined portions in its periphery and a clamping means capable of clamping the front end of the paper, wherein the clamping mean includes at least a delivery roller 32 parallel to the drum cylinder, the delivery roller 32 clamps a predetermined portion of the printed paper and drives the printed paper to rotate with the drum cylinder. In this way, the printed papers are accumulated and clamped on the delivery roller 32. Also, the delivery roller 32 may discharge several pieces of printed papers out of the drum cylinder 20. This device, however, has the following disadvantages.\nFirstly, the device needs to employ a clamping means to clamp the front edge of the paper by its pressure, such that the accumulating device can then operate normally. In this way, since it is necessary for the clamping means to open and close repeatedly, its clamping force will decrease gradually due to spring fatigue, which causes clamping failure, and thus the accumulating device fails to operate normally.\nSecondly, as is shown in FIG. 1B, it is required for the device to rotate positively and negatively by means of two motors, i.e., with the reversible rotating of a first motor 66, the device controls the drum cylinder 20 to rotate reversibly and thus implements paper accumulating; with the reversible rotating of a second motor 82, the device controls the delivery roller 32 and pressing roller 36 to rotate positively and negatively and thus implement paper discharging. The structure involved is complicated and the cost for electric control is relatively high."}
{"text": "The present invention relates to methods and apparatuses permitting commands to be issued to a computer program presenting a virtual reality environment to users of the program through necessary hardware interfacing with the computer program.\nInteracting with a computer using a virtual reality environment requires means to give commands to the computer. Standard methods, such as use of keyboard and/or mouse plus a pull-down menu, of giving commands to a computer application do not function well in, and are not suited for, a virtual reality system. Voice commands are useful, but they are at present limited to only simple suites of commands because of the heavy computational power required for voice word recognition. Voice commands also require the user to memorize the exact phrasing of each command.\nThe present invention allows the programmer to easily create a collection of virtual buttons and menus in a virtual reality environment and the user to select commands by looking at a virtual button.\nPerhaps the closest art is provided by the \"haptic button\", wherein the user inserts his or her hand in a device which simulates the tactile sensations which would be caused if the hand contacted a real object corresponding to an object in the virtual reality environment. Buttons can be simulated thereby and depressed by tactile push. The problem with this means is that it is also computationally expensive and requires an additional device which is expensive and prevents both hands from being free for manipulation of objects outside the virtual reality environment."}
{"text": "Bimetallic tubulars with highly corrosion resistant inner liners can be utilized for reducing drill pipe corrosion in super sour well environments containing elemental sulfur. However, when the bimetallic tubulars are joined, alloying of the two distinct materials comprising the bimetallic tubular may occur and significantly reduce the bimetallic tubular's corrosion resistance and strength at the point of the weld. What is needed in the art is a method for welding and a weld configuration which preserves the continuity of the corrosion resistant inner liners of bimetallic tubulars and which avoids alloy formation of the different materials of the tubular layers at the weld joint.\nThe welding method and weld configuration of the instant invention advantageously avoid alloying of the discreet materials comprising the layers of the bimetallic tubulars when the weld joints are formed. Thus the possibility of decreased corrosion resistance and strength resulting from alloying of the two materials composing the bimetallic tubular in the weld configuration is eliminated."}
{"text": "The present invention generally relates to contactors, and especially to a contactor making a contact to electronic components such as a large-scale integrated circuit and a contact process that uses such a contactor.\nIn recent years, the development of production technology in the field of semiconductor substrate is remarkable. Associated with this, interconnection patterns of large-scale integrated circuits are miniaturized, and along with this, the terminals of large-scale integrated circuits are also miniaturized. Further, the number of the terminals used in an LSI is increasing with remarkable rate.\nThe demand of miniaturization and high-density mounting is acute especially in the apparatuses that use a large-scale integrated circuit. For example, the number of mobile apparatuses (cellular phones, mobile personal computers, video integrated cameras etc.), in which downsizing is demanded, or high-performance computers in which the distance between adjacent LSIs has to be minimized for guaranteeing high speed operation, is increasing rapidly.\nThe foregoing demand also affects on the shipment mode of LSIs. Thus, the cases are increasing for shipping unpackaged LSI chips while guaranteeing the operability thereof as known as KGD (Known Good Die), or shipping the LSIs in the form of CSP (Chip Size Package), which is a small-sized package having the size of a chip.\nFrom these circumstances, there is a need of a contactor capable of making a contact with a large number of miniature pin terminals with certainty for testing the LSIs.\nAlso, from the viewpoint of efficient test of LSIs, there is emerging a need of full test, in which all the tests such as final test (FT) or burn-in (BI) test are conducted for each of the LSIs in the state of wafer, before the wafer is divided into individual LSI chips. By using the full test in the state of wafer, following effects are expected.\nFirst, the efficiency of handling is improved as compared with the case of conducting the testing on separate chips. When the size of the chips is different, it should be noted that the compatibility of the handling equipment used for the testing is lost. In the case the testing is conducted in the state of wafer, on the other hand, it becomes possible to convey the wafers of standard outer size one after another. Further, it becomes possible to control the defect information in the form of wafer map.\nIn the case the art of wafer-level CSP, which is subjected to intensive research and development in recent years, is used, it is potentially possible to conduct the entire process steps from the wafer process up to the assembling step in the form of the wafer. Thus, if it becomes possible to realize the test process in the state of wafer, the entire steps from the wafer process to the packaging process (assembling process) can be conducted in the state of the wafer, and the efficiency of production of LSI chips would be improved significantly.\nHowever, as noted before, it has been difficult to realize a conductor that can contact the terminals of plural LSIs on a wafer or the terminals of entire LSIs on the wafer, in view of miniaturization of individual LSIs and hence the terminals thereof, and further in view of increasing number of the terminals.\nHereinafter, typical examples of conventional contactors (probe cards or sockets) will be summarized.\n(1) Needle Type Mechanical Probe\nA needle type mechanical probe has a construction of disposing a plurality of needles formed of a tungsten wire and the like on a contactor substrate of an insulating material in correspondence to respective terminals of the LSIs to be tested. Generally, a cantilever structure has been used in which the needles are provided so as to extend obliquely over the LSI wafer. Further, there is a proposal of disposing needles in a vertical direction to the terminals of the tested LSIs, by providing resiliency to the needles.\n(2) Membrane Type Probe\nA membrane type probe has a structure of a film circuit having a metal projection (referred to hereinafter as “bump”) for the contact electrode of the probe.\n(3) Anisotropic Conductive Rubber\nAn anisotropic conductive rubber uses an elastic rubber as an insulating base material and has a structure in which a conductive material that extends only in the thickness direction of the rubber base material such as a metal wire is incorporated.\nFurther, the Japanese Laid-Open Patent Application 10-111316 official gazette discloses a contactor in which an end part of the wire, extending out from an edge of a substrate, is used for a connection terminal that makes a contact with the semiconductor device to be tested and the terminals extending at the other edge of the substrate is used for the measurement terminal.\nOn the other hand, the abovementioned needle type mechanical probe has problems such as:\na) High cost of forming a large number of needles or pins individually;\nb) Limitation in the precision of the needle tip due to the construction in which a large number of needles are arranged individually; and\nc) Restriction imposed on the arrangement of the needles in view of the oblique arrangement of the needles.\nThus, it has been difficult to use a needle type mechanical probe for the contactor that makes a contact with plural LSIs at the same time.\nFurther, the abovementioned membrane type probe has the problems as summarized below.\na) Individual contact electrodes cannot move freely\nBecause the contact electrode is connected to a polyimide layer forming the insulating substrate, the movable range of the individual electrode is limited. Also, in view of the fact that the contact electrode is formed of a metal bump of a hard metal that lacks flexibility, there arises a problem in that a defective contact may be caused in the case there exists a change of height between the bump electrodes.\nb) High Cost\nThe bump constituting a contact electrode is generally formed by plating a metal. Thus, it takes time for forming a bump and the cost is inevitably increased.\nFurthermore, the anisotropic conductive rubber has problems summarized as below.\na) Limited Lifetime\nIn the case it is used at high temperatures, in particular (it should be noted that the BI test is carried out usually at the temperature of 125° C. or more), the rubber part undergoes plastic deformation, and it can be used for only a dozens of time at best.\nb) It cannot be used for the case of narrow electrode pitch.\nBecause it is difficult to incorporate a conductive material into a rubber, the pitch of 200-150 μm is thought as being a practical limit.\nAlso, the contactor disclosed in the Japanese Laid-Open Patent Application 10-111316 official gazette has a disadvantage in that it is difficult to achieve a contact for all the electrode terminals in the case there exists a variation of height in the electrode terminals of the semiconductor device to be tested, in view of insufficient stroke of the contactor terminal in the longitudinal direction. Further, there can be a possibility that sufficient contact cannot be achieved in the case the thermal expansion coefficient of the semiconductor device to be tested is different from that of the base material of the contactor as a result of displacement caused at the time of the high temperature test such as in the case of the burn-in test.\nThus, the conventional contactors have the problems summarized as follows.\n1) Insufficient stroke of the contactor contact terminals\nThus, in such a conventional contactor having insufficient stroke (elastic deformation) for the contact part, a sufficient contact is achieved only when it is caused to make a contact with an aluminum pad on a wafer, in which the variation of the height is relatively small, or in the case it is used for a narrow area of the wafer. However, in the case of a wafer level CSP or molded packages formed by a simultaneous molding process, the terminals or balls on the package generally have a large variation of height, and it becomes difficult to achieve a sufficient contact by using such a conventional contactor having a small stroke or elastic deformation in the contact part. This problem becomes especially serious in the case the contactor is used to make a simultaneous contact with a thin wafer, which tends to have a warp. It should be noted that the warp of a wafer tends to increase when the thickness of the chip is reduced.\n2) Positional deviation by the difference of the thermal expansion coefficient\nThe LSI wafer to be tested is generally formed of silicon (Si). It should be noted that the coefficient of linear thermal expansion is about 3 ppm in the case of silicon, while in the case of an insulated substrate used for the contactor such as a resin, the value of the coefficient of linear thermal expansion becomes several ten ppm (13˜30 ppm, for example). Thus, even in the case the contactor is contacting properly at ordinary temperatures, the location of the contactor may be deviated due to the difference of the coefficient of linear thermal expansion when it is used at a high temperature as in the case of a BI test. In an extreme case, the contactor may miss the intended terminal or make a contact with a terminal next to the intended terminal. Further, a similar problem occurs also in the case of the wafer level CSP or packages molded by a simultaneous molding process, in view of the fact that the coefficient of linear thermal expansion is different between the insulated substrate material and the package material such as seal resin. In the case polyimide is used for the insulating substrate, for example, the linear thermal expansion coefficient takes a value of about 13 ppm and there can occur a displacement as much as 100 μm at the peripheral part of an 8-inch wafer, which has a diameter of about 100 mm, when it is heated to 125° C., even in the case the contactor is aligned properly at ordinary temperature."}
{"text": "This patent specification refers to, and hereby incorporates by reference, the items listed at the end and identified by reference (ref) numbers in parenthesis in the discussion below.\nPhase contrast imaging is a magnetic resonance imaging (MRI) technique that can be used to visualize moving fluid. It is typically used for MR angiolgraphy that does not require introducing a contrast agent in the subject. Phase contrast imaging relies on the assumption that spins that are moving in the same direction as a magnetic field gradient develop a phase shift that that is proportional to the velocity of the spins. In principle, bipolar gradients (two gradients with equal magnitude but opposite direction) are used to encode the velocity of the spins. Stationary spins undergo no net change in phase after the two gradients are applied, but moving spins experience a different magnitude in phase response to the second gradient compared with the first. The detected net phase shift can be used to calculate the velocity of the spins, or an image obtained with a bipolar gradient can be subtracted from an image of the same slice obtained without using a bipolar gradient to thereby determine velocity from net phase shift and highlight flow. Gradient moment nulling can be used to balance the phase of all gradient pulses to achieve only velocity phase shifts in the echo signal from an effective bipolar pulse.\nThe most common and available method for hemodynamic MR velocity imaging on commercial clinical scanners is believed to be two-dimensional (2D) phase contrast (PC) imaging using a bipolar gradient pulse (ref 1) to encode the signal phase with velocity in cine gradient echo (GRE) sequences with low flip angle radio frequency (RF) pulses, and cardiac gating (cine-PC) through the cardiac cycle (ref 1, ref 2). Cine-PC 2D imaging measures flow velocity in a single slice plane and scan time is increased proportional to the number of repeated slice planes through arteries or cerebrospinal fluid (CSF) passageways."}
{"text": "Semiconductor devices and integrated circuits are typically manufactured on a single semiconductor wafer. The dies of the wafer may be processed and packaged with other semiconductor devices (e.g. antenna) or dies at the wafer level, and various technologies have been developed for the wafer level packaging."}
{"text": "The present invention relates generally to an apparatus for measuring the amount of skew or sinusoidal variation in the edges of moving plates or sheet material. It is particularly intended for use in preparing steel sheet and strip for the manufacture of tin plate.\nIn the manufacture of tin plate from sheet and strip it is normal practice to subject a slab of steel to a hot rolling operation wherein rolling at the last finishing stand is conducted above the upper critical temperature of the metal. In such a hot rolling operation the strip is wound into a coil at temperatures ranging from 1050.degree. to 1200.degree. F., depending on the end use of the product and its desired metallurgical characteristics. The hot roll strip is then subjected to a continuous pickling operation wherein the strip first passes through cold working equipment which fragments surface scale and facilitates acid attack prior to actually passing through the pickling solutions at uniform speed to complete oxide removal, and is thereafter followed by cold water spray rinses and, if necessary, neutralizing alkaline solutions. The strip is then recoiled.\nFollowing pickling, cold reduction of the coil takes place. The reduction in thickness of the strip may be as great as 90% and is carried out at exceedingly high rates of speed on the order of 1 mile per minute. Care must be exercised to obtain flat strip from the cold reduction mill in order to secure uniform results further on in the tinning process. Strip shape can be distorted by edge wrinkles or center buckles, most of which is caused by excessive uneven rolling pressures. Such distortion results in noncylindrical out of round coils of metal.\nPrior to annealing the coils received from the coil reduction operation are uncoiled and subjected to electrolytic cleaning, rinsing and air drying and are thereafter recoiled. The strip may be annealed on a continuous line, which once again necessitates uncoiling and recoiling the strip at relatively high temperatures and line speeds.\nThe annealed strip may then be subjected to temper-rolling to secure the desired hardness and surface texture, to impart the required mechanical properties to the product made from the strip and to effect final flattening of the strip. Again, excessive rolling at the edge or center can cause edge or center fullness.\nIn the course of the foregoing operations the strip is coiled and uncoiled several times. Apart from any edge or center defects, improper recoiling can cause noncylindrical coils with resultant defects and reduced acceptance of the ultimate tin plate product.\nA typical coil when treated in the foregoing operations may weigh from 30,000 to 40,000 pounds. Its length may be approximately 25,000 feet. It can be readily appreciated that handling such large units of product, with frequent coiling and uncoiling, can result in recovery of coils that are noncylindrical and undesirable.\nIn tin plate manufacture a coil preparation line is thus used for edge shearing and inspection of the strip at a location prior to the tin coating operation. As in prior operations, the coils of steel are placed on a motor driven mandrel and uncoiled through tension bridle rolls, side shears and then recoiled again on a motor driven mandrel while moving through the line at very high speeds such as 2400 fpm.\nNoncylindrical coils when processed through the rotary side shear can result in edges on the strip which are not straight but instead have sinusodial variations. Such variations produce inferior product and can result in rejection of products and great loss to the manufacturer. Accordingly, installation of a measuring device on the coil preparation line to detect edge variations as the coil is sheared can be very advantageous and can provide information enabling corrective action to be taken before inferior products are produced. One useful arrangement is to locate the measuring system on the exit side of the side shears.\nThe lack of symmetry of an out of round coil and the sinusodial oscillations of the steel strip which result when it is unwound from the coil preparation line mandrel as it is fed into side trimmer knives are definite disadvantages. The amplitude of the oscillations tend to increase with the rotational speed of the mandrel and the strip speed. Such oscillations produce an alternating tension from side to side on the strip as it is fed into the shear, thus producing a sinusodial pattern of lateral weave cut into the sides of the strip by rotating shear knives. The period length of the sinusodial variation is dependent on the circumference of the payoff coil. When the edge of the strip deviates from a straight line over a given longitudinal length, that length of strip becomes unusable for prime tin plate products.\nThe purpose of this invention is to determine dynamically the maximum amplitude of the deviation from a straight line for a given wave length of strip, e.g. 0.025 inches in 5 feet. This is accomplished by taking three separate lateral measurements along the edge of the strip along with a measurement of wave length and thereupon calculating maximum amplitude of the lateral variations."}
{"text": "Field\nThe subject application relates to data transmission, e.g., asymmetrical data transmission over simplex and duplex channels.\nBackground\nCellular telephones are tremendously popular. It is estimated that at the end of 2007 the total worldwide subscriber rate reached 3.3 billion. Close to 80% of the world's population enjoys mobile telephone coverage, a figure that will only increase. As cellular telephones gain popularity, their functionality has increased also. Standard service includes voice calling, caller ID, call waiting, and voice mail. Service providers also offer text messaging, push mail, navigation, and even a high-speed internet connection directly to the telephone through the use of protocols such as those included in High Speed Packet Access (HSPA).\nHSPA is a collection of wireless protocols that improve upon the performance of existing Universal Mobile Telecommunications System (UMTS) protocols. High-Speed Downlink Packet Access (HSDPA), a standard within HSPA, increases data packet transfer performance by using improved modulation schemes. These improved schemes better utilize existing radio bandwidth provided by UMTS. HSDPA currently supports downlink speeds of 1.8, 3.6, 7.2, and 14.4 Mbit/s. Long Term Evolution (LTE) is a promising standard for the next generation (4G) of mobile broadband networking.\nMultiple-input and multiple-output (MIMO) requires the use of multiple antennas at both the transmitter and receiver. The signals from the antennas are combined to minimize errors and optimize data speed, providing better range and performance. However, the use of multiple inputs and outputs requires a device to utilize the same radio spectrum frequency. The United States presently uses the GSM-850 and GSM-1900 radio spectrum frequencies for cellular transmissions. GSM-850 uses 824-849 MHz for uplink and 869-894 MHz for downlink, providing channel numbers 128-251. GSM-1900 uses 1850-1910 MHz to uplink and 1930-1990 MHz to downlink, providing channel numbers 512-818. The MIMO concept defined in Third Generation Partnership Project Revision 7 (3GPP R7) and Revision 8 (MIMO R8), incorporated by reference herein in their entirety into this disclosure, requires the use of the same radio spectrum frequency for both transmission paths. These frequencies and antennas are used in spatial multiplexing or transmission diversity mode according to radio conditions. This allows for multiple simultaneous data streams, thereby increasing the data transmission rate.\nMIMO R8 also requires twice the amount of antennas at both the transmitter and the receiver locations, even though the transmission takes place across a single frequency band. This creates interference in the signal, which decreases the actual gain in bandwidth created by MIMO R8. The additional signal used in MIMO R8 is another two-way transmit path. Although MIMO R8 can have up to four transmit paths, the uplink bandwidth is still equivalent to the downlink bandwidth, because each additional transmit path adds a duplex channel.\nDemanding data services for individual users can exceed the capabilities of a single frequency carrier and/or radio path for a variety of transmission technologies. In this case, the capacity of multiple bi-directional frequency carriers and/or radio paths are combined, or “bonded” for the single demanding user. Multiple pre-existing bi-directional transmission pairs are allocated to the demanding user and traffic is spread across them. The Federal Communications Commission (FCC) recently auctioned the 700 MHz frequency spectrum. AWS-700 uses 776-794 MHz for uplink and 746-764 MHz for downlink.\nAs is, these transmission techniques offer useful means to boost individual peak throughput within the capabilities of the available transmission technology. However, bi-directional frequency carriers and/or radio paths, and the equipment required to use them, are best utilized if the data load and equipment capabilities are symmetrical. Unfortunately this is often not the case. Traffic for most data, audio, and video applications is heavily weighted in the downlink, server to user, direction. Roughly eight times as much data is downloaded to as is uploaded from mobile devices. The number of duplex signals available may limit these downlink requests. Many frequency bands are not intended for and not licensed for transmission by an individual user. Currently there are many of these frequencies available for downlink only which are being underutilized. Subscriber equipment, especially wireless, is also limited in terms of available space, power (battery life for mobile devices) and cost. The need for subscriber equipment to simultaneously transmit on all bonded frequency carriers and/or radio paths is therefore an unnecessary burden from an equipment complexity, cost and power perspective.\nWhat is needed is a system that utilizes a downlink only channel to supplement the bandwidth of a conventional duplex channel to distribute the data load."}
{"text": "1. Field of the Invention\nThe present invention relates to cephalosporin antibiotics. More particularly, the present invention includes novel (7R)-7-(acylamino)-3-(arylthio)-3-cephem-4-carboxylic acids and their pharmacologically acceptable salts and prodrugs, their methods of production and use. These compounds exhibit antibiotic activity against a wide spectrum of organisms, including organisms which are resistant to .beta.-lactam antibiotics.\n2. Review of the Background Art\nOver the past three decades a large variety of antibiotics has become available for clinical use. One class of antibiotics which has seen remarkable growth are the cephalosporins (shown generically below), over 70 of which have entered clinical use for the treatment of bacterial infections in mammals since 1965. The cephalosporins exhibit their antibacterial activity by inhibiting bacterial peptidoglycan biosynthesis, and have been extremely effective in treating a wide variety of bacterial infections. Cephalosporins that are said to have antibacterial activity are described in U.S. Pat. No. 3,992,377 and U.S. Pat. No. 4,256,739. ##STR1##\nUnfortunately, the wide-spread and indiscriminant use of these antibiotics has led to a rapid increase in the number of bacterial strains which are resistant to these compounds. Most importantly, this resistance has emerged among clinically important microorganisms which threaten to limit the utility of presently available cephalosporin antibiotics. In particular, resistant strains of Salmonella, S. pneumoni.ae butted., Enterobacteriac.ae butted., and Pseudomonas have emerged which threaten to undo many of the strides made in reducing mortality and morbidity from bacterial infections.\nBacterial resistance to cephalosporins follows three major pathways: a) the development of .beta.-lactamases capable of inactivating the .beta.-lactam ring of the cephalosporin; b) decreased cephalosporin penetration into the bacteria due to changes in bacterial cell wall composition; and c) poor binding to penicillin-binding proteins (PBPs). The latter pathway is especially important, as the binding of .beta.-lactams to PBPs is essential for inhibiting bacterial cellwall biosynthesis. Certain Gram-positive bacteria, namely methicillin-resistant Staphylococcus aureus (\"MRSA\") and enterococci are highly resistant to .beta.-lactam antibiotics. Resistance in MRSA is due to the presence of high levels of an unusual PBP, PBP2a, which is insensitive, or binds poorly, to .beta.-lactam antibiotics. The activity of .beta.-lactam antibiotics against PBP2a-containing organisms has been shown to correlate well with the binding affinity of the antibiotic to PBP2a. Currently, the glycopeptides vancomycin and teicoplanin are primarily used for MRSA bacteremia. The quinolone antibacterials and some carbapenems, such as imipenem, have been reported to be active against a few MRSA strains, but their use is restricted due to emerging resistant MRSA strains.\nExperimental compounds which may possess utility as anti-MRSA or anti-enterococcal bactericides include the glycylcyclines (see, e.q., P.-E. Sum et al., J. Med. Chem., 37, (1994)), FK-037 (see, e.q., H. Ohki et al., J. Antibiotics, 46:359-361 (1993)), RP-59,500 (see, e.g., S. K. Spangler et al., Antimicro. Agents Chemother., 36:856-9 (1992)), the everninomycin complex (see, e.g., W. E. Sanders et al., Antimicro. Agents Chemother., 6:232-8 (1974)), the 2-(biaryl)carbapenems (see, e.g., U.S. Pat. No. 5,025,006), 3-(benzothiazolylthio)cephems (see, e.g., EP Application No. 527686), 3-(thiazolylthio)carbacephems (see, e.g., R. J. Ternansky et al., J. Med. Chem., 36:1971 (1993) and U.S. Pat. No. 5,077,287) and arbekacin (S. Kondo, et al. J. Antibiotics 46:531 (1993)."}
{"text": "Japanese non-examined laid-open Patent Publication JP H9-37670 A discloses an absorbent panel for pets. The absorbent panel has a liquid-absorbent core panel between a top sheet and a back sheet. The core panel has a relatively thin central region and a relatively thick edge region. Thus, the absorbent panel is configured such that an edge region of the top side of the absorbent panel is higher than its central region or protrudes on the top side, while its back side is flat. With such a structure, urine is prevented from leaking out from a peripheral edge region of the absorbent panel."}
{"text": "1. Field of the Invention\nThe present invention relates to a defect detector and a method of detecting a defect and more particularly, to a defect detector for inspecting a surface of an electrophotographic photoreceptor used for copiers, printers, and facsimiles.\n2. Discussion of Background\nIn processes of manufacturing electrophotographic photoreceptors used for copiers, printers, and facsimiles, the photoreceptors occasionally have defective surfaces such as scratches, foreign materials, and coating irregularities. As a method of detecting the defects, a visual inspection is typically performed by an inspector. However, since the visual inspection has fluctuations due to individual differences among inspectors, various automatic detecting methods and apparatus have been suggested.\nFor example, as a prior art apparatus and method, Japanese Laid-Open Publication No. 9-325120 (hereinafter referred to as a Prior Example 1) discloses a detector detecting the following defects by irradiating an electrophotographic photoreceptor from a light source, receiving light reflected from the photoreceptor with a camera, and producing an electrical signal representing the difference between:                (1) a concavo-convex defect having a high rate of change of surface concavity and convexity, such as pin holes, dents, coating scratches, and dusts; and        (2) a blocky defect having a low rate of change of surface concavity and convexity, such as coating irregularities of the photosensitive layers.        \nIn addition, Japanese Laid-Open Publications Nos. 7-128240(hereinafter referred to as a Prior Example 2) and 7-128241(hereinafter referred to as a Prior Example 3) disclose a detector taking plural photographs of a concavo-convex defect and a blocky defect with a light source and plural charge coupled device (hereinafter “CCD”) cameras or a CCD camera and plural light sources.\nHowever, in an optical system of the Prior Example 1, both the concavo-convex defect and the blocky defect are photographed without distinction, and a decision to pass or fail has to be made in accordance with the defect having a severer acceptance criterion. Therefore, even a normally good photoreceptor is disposed of as a defective one, or a photoreceptor, which the defect detector judges to be rejected, has to be inspected again by an inspector. In addition, there is an advantage that the optical system is easily controlled because of its few controllers, but there is a disadvantage that when sensitivity of photographing either the concavo-convex defect or the blocky defect is increased, sensitivity of photographing the other defect is decreased.\nFurther, in the Prior Examples 2 and 3 , two light sources or two CCD cameras can increase sensitivities of photographing both the concavo-convex defect and the blocky defect. However, inspection time is long because of its plural times of photographing. In addition, the control operation is complicated because plural CCD cameras are combined with a light source or plural light sources are combined with a CCD camera.\nBecause of these reasons, a need exists for a defect detector which can efficiently detect a concavo-convex defect and a blocky defect by photographing a defect having a high rate of change of surface concavity and convexity and a defect having a low rate of change thereof with separate optical systems, and which can reduce losses due to excessive inspections."}
{"text": "It is sometimes necessary or otherwise desirable to identify and treat a bodily passage, such as a sinus cavity. For example, when sinus cavities become blocked, balloon sinuplasty—the dilation of the sinus using a balloon catheter—provides an alternative to radical surgical approaches to unblocking the sinus cavity. Conventional procedures utilize two separate devices to accomplish balloon sinuplasty, a rigid scope and a guiding catheter. The scope is utilized to locate the point of treatment and to provide visualization throughout the treatment process and the guiding catheter is used to guide a treatment device to the point of treatment.\nAdvancing a scope and guiding catheter separately towards a point of treatment has significant drawbacks because it increases the complexity of the procedure and the amount of time required to complete the procedure. For example, the individual performing the procedure must use both hands to manipulate the scope and guiding catheter. Doing so is awkward and, in some cases, requires an assistant to complete the procedure. Furthermore, due to the tortuous anatomy of many bodily passages, physicians generally must maintain a variety of scopes and guiding catheters that have many different angle configurations so that particular devices with appropriate configurations can be selected at the time of treatment.\nScopes have been developed that include a deflectable tip, which allows physicians to view various portions of a bodily passage during the performance of a procedure. However, these scopes still present significant drawbacks. For example, scopes that include deflectable tips are rigid in nature, which limits their ability to be navigated through the tortuous anatomy of a bodily passage. In addition, these scopes do not incorporate a flexible material or a mechanism for transitioning between a rigid configuration and a flexible configuration during the performance of a procedure. This prevents adjustment of the scope during the procedure such that the scope can be adapted to the anatomy of a bodily passage, for example, when it is desired to pass the scope through an opening defined by the bodily passage.\nA need exists, therefore, for improved medical devices and methods for identifying and treating a bodily passage, such as a sinus cavity."}
{"text": "Carbon dioxide exists as a low-density gas at standard temperature and pressure conditions and possesses phase boundaries with a triple point (Solid-Liquid-Gas co-exist in equilibrium like a glass of ice cubes and water) and a critical point (Liquid-Gas have identical molar volumes). Through pressure or temperature modification, carbon dioxide can be compressed into a dense gas state. The term ‘Dense Phase Carbon Dioxide’ is used herein to describe all phases of carbon dioxide: liquid state, supercritical state, dense gas state, and solid-state. These states have densities that are within the range of liquid-like or near-liquid substances.\nCompressing carbon dioxide at a temperature below its critical temperature (C.T.) liquefies the gas at approximately 70 atm. Cooling liquid-state or gas-state carbon dioxide to its freezing point causes a phase transition into solid-state carbon dioxide. Compressing carbon dioxide at or above its critical temperature and critical pressure (C.P.) also increases its density to a liquid-like state, however there is a significant difference between compression below and above the critical point.\nCompressing carbon dioxide above its critical point does not effect a phase change. In fact, carbon dioxide at a temperature at or above 305 K (88 F) cannot be liquefied at any pressure, yet the density for the gas may be liquid-like. At the critical point the density is approximately 0.47 g/ml. At or above this point carbon dioxide is termed a supercritical fluid (SCF). Supercritical carbon dioxide can be compressed to a range of liquid-like densities, yet it will retain the diffusivity of a gas. Continued compression of supercritical carbon dioxide causes continued increase in density, approaching that of its liquid phase.\nSolid-state carbon dioxide is useful for removing particulates and trace organic residues from surfaces, using a process typically called Snow Cleaning or Snow Departiculation. Similar to liquid and supercritical fluid cleaning agents, solid-state carbon dioxide's cleaning power can also be described in both physical and chemical terms. The process of snow departiculation can be described as a kinetic energy transfer process called “Linear Momentum Transfer” in accordance with the following vector quantity:\nP=MV, where\nP—Linear Momentum of Solid Carbon Dioxide Particle or Surface Particle\nM—Mass of Solid Carbon Dioxide Particle or Surface Particle\nV—Velocity of Solid Carbon Dioxide Particle or Surface Particle\nA stream of solid carbon dioxide particles having significant mass and velocity impact a stationary surface particle causing the surface particle with a given mass to accelerate away from the surface to a given velocity in accordance with the following equation:Vsp=(Mcp/Msp)Vcp, where\nVsp—Velocity of Surface Particle\nMcp—Mass of Carbon Dioxide Particle\nMsp—Mass of Surface Particle\nVcp—Velocity of Carbon Dioxide Particle\nThe physical energy transferred during a snow departiculation process is usually sufficient to overcome strong electrostatic and intermolecular adhesive forces, commonly referred to as Van der Waal's forces, that hold small particles to the surface.\nThe mechanism for the removal of trace organic films using snow is not fully understood, but has been postulated to be a combination of momentum transfer and a phase change of minute solid carbon dioxide particles from solid-state to liquid-state (compression) and subsequent solutioning of trace surface residues. According to the phase diagram for carbon dioxide, a minimum impact compression of approximately 6 atm (88 psi) at 195 K is required to produce a liquid interphase. Energy transformations are possible other than the formation of a liquid phase, including particle fragmentation or shearing, gas phase transition (sublimation), and temperature rise in the solid (thermal energy) at impact.\nSolid carbon dioxide is being used in a number of commercial product cleaning applications to remove trace organic and inorganic residues and particulates. Liquid carbon dioxide is rapidly expanded (joule-thompson expansion) through an orifice of a valve to form a mixture of subcooled gas state and solid state carbon dioxide—referred to as “snow” or “dry ice”.\nSolid carbon dioxide is applied in conventional applicators according to two types of applicators, described as Type I and Type II snow cleaning applicators as follows:\nType I Snow Applicator: Liquid carbon dioxide, stored in a high pressure bottle, is expanded from 850 psi at 298 K through a suitable nozzle into gas state (the propellant) and solid state carbon dioxide (the cleaning agent) and directed at a substrate. Conventional Type I applicators are commonly used in precision cleaning applications at close proximity to a substrate and have relatively simple operation and low-cost designs.\nType II Snow Applicator: Liquid carbon dioxide is first expanded into solid carbon dioxide using a suitable “dry ice machine”, packed into dry pellets of uniform size, or shaved into a powder, and then fed into a spray apparatus using compressed air to propel the solid carbon dioxide from a spray nozzle. The air and solid mixture impacts the surface. A Type II applicator is typically used for cleaning large rigid structures because of its more aggressive action at close proximity (i.e., for coating removal) and long-range particle cleaning action (large and hard snow pellets). However, Type II equipment and operational costs are significantly higher than Type I systems.\nType I and Type II applicators include:\n1. Fixed position applicators\n2. Pistol grip applicators\nConventional applicators are designed to have a single spray pattern, with various interchangeable nozzle designs for different substrates and surface cleaning applications. Type II designs can also vary impact energy through control of compressed air pressure whereas Type I designs cannot. Disadvantages associated with these mechanical designs include:                1. Fixed spray pattern.        2. Non-interchangeability of applicator designs (fixed<−>handheld<−>robotic).        3. Bulky configurations.        4. Uncontrolled tribocharging (electrostatic buildup) of non-metallic substrates such as plastics.        5. Rapid localized substrate cooling and subsequent deposition of contaminating residues.        6. Ineffective deep hole cleaning.        7. Expensive equipment costs.        \nConventional snow cleaning applicators (Type I and II) suffer from the following disadvantages:                1. Impact energy and the amount of snow particles available at the surface decreases as the distance from the expansion valve to the applicator nozzle increases. Type I applicators must have an expansion valve located close to the nozzle and the nozzle must also be very close to the substrate to effectively remove residues.        2. Entrainment of ambient air which often contains moisture, particles, and other contaminating residues which condense onto surfaces which have been supercooled by the snow particle/gas stream.        3. Externally applied environmental control measures such as heated air and particle control hinder the cleaning performance and are applied so generally that localized condensation, particle entrapment, or tribocharging still occur. Conventional applications employ macro-environmental control (clean rooms, infrared heaters etc.) measures. Type I snow applicators cannot be used in relatively uncontrolled environments.        4. The process of expanding liquid carbon dioxide into solid state and subsequent contact of solid state carbon dioxide with surfaces causes a phenomenon called tribocharging, whereas the solid carbon dioxide (primary dielectric) builds electrostatic charge of up to 5 to 15 mJ at 10 KV to 20 KV as it contacts a substrate (secondary dielectric). This type of electrical charge build-up can be extremely damaging to microelectromechanical devices (or can induce latent ESD defects) and will cause a departiculated surface to become an attractor (magnets) of airborne particles following snow cleaning operations. Electrostatic effects can be caused through direct contact of charged solid carbon dioxide particles with the substrate which causes a discharge event or current flow through the surface (direct discharge) or may be caused through electrostatic field exposure and subsequent charging of the surface (induced charging).        5. In many applications, spray cleaning is performed independent of and prior to operations such as microwelding, adhesive bonding and thermal curing and soldering. Moreover, following production operations such as CMP the substrate is wet with aqueous residues and must be dried prior to snow cleaning operations. A method and apparatus is needed to serve as an integrated simultaneous drying, cleaning and production tool.        6. Xenon flashlamp technology is used with solid carbon dioxide (Type II) to remove old paint from aircraft surfaces. This type of technology uses an intense UV radiation (not a laser) burst to pyrolize substrates which produces a large amount of heating radiation as a by-product. The pyrolized paint is swept away from the substrate using a flow of carbon dioxide pellets. This technology is large and bulky and cannot be used to precisely clean small parts commonly found in the semiconductor, electrooptical and electronics markets. A precision coherent photon-based technique is needed to remove small contaminants from intricate assemblies.        7. In some applications, solid phase carbon dioxide chemistry requires physicochemical modification to provide enhanced separation and surface finishing capabilities. To date, no effective technique has been demonstrated to accomplish this requirement.        8. Type II applicators also suffer from being too aggressive (i.e., substrate damage), very noisy, bulky and too costly for most precision substrate cleaning applications.        9. Conventional snow cleaning applicators do not lend themselves to integration with production processes such as stamping, welding, bonding, curing and abrasive surface finishing operations because of the aforementioned problems discussed above.        \nConventional snow cleaning processes do not have a method for real-time analysis of cleaned surfaces to accept or reject a particular cleaning operation. This is especially advantageous for in-line continuous quality control monitoring of surface cleaning performance.\nConventional ESD Control Methods used with Solid Carbon Dioxide:\nAir Ionization—air is ionized using a DC or AC ionizer that is then flushed over an affected surface. The problem with this approach is that flowing air induces contamination through introduction of humidified air and potential particles. Also ionizing air impingement requires flooding the surfaces to be cleaned. This process can subtract from the cleaning energy. Moreover, the charges present within the structure of the cleaning agent are not reduced effectively using this technique.\nU.S. Pat. No. 5,409,418 proposes a nozzle-mounted secondary gas ionizer which surrounds the snow stream with oppositely charged ions during impingement. U.S. Pat. No. 5,725,154 proposes neutralizing charges during snow cleaning following each cleaning pulse with a separate propellant gas neutralization pulse.\nMost prior art suffers from these typical drawbacks:                Impossible to precisely control charges being delivered to a substrate—each substrate and atmosphere is different.        The portion of a substrate being impacted by the sublimable cleaning agent is not affected by neutralizing ions—only the circumference of the snow spray is affected.        Backside or nearby electrostatic charging due to electric fields is not affected by these techniques—electric fields pervade the materials creating complexly charged surfaces.        \nNuclear Ionization—the substrate is exposed to radioactive particles (alpha). This process is line-of-sight and very short range. Obstructions of the smallest variety will eliminate beneficial ionization using this technique.\nFong '786, referenced herein, uses nuclear ionization to reduce accumulated electrostatic charges contained on solid carbon dioxide stored and mixed within a storage hopper and prior to and during delivery into a high pressure feed line. Fong '786 suffers from all of the drawbacks cited above.\nGrounding—the substrate in grounded to earth using a suitable resistor to bleed charges at an acceptable rate. The main problem with this approach is that the electrostatic charge and electrical overstress are not effectively controlled on non-conductive substrates.\nAntistatic Chemicals—this approach is the most effective on preventing charge creation by the cleaning agent. However this method tends to, by itself, become a source of chemical contamination within the cleaning process. To date no use of antistats within cryogenic cleaning agents is known\nMoreover, in cleaning quartz lenses, as well as many other non-conductive substrates it is difficult to control electrostatic charging of the quartz substrate during sublimable spray cleaning. Flooding the surfaces with ionized air only works prior to and following snow cleaning. During snow cleaning, as much as 2000 volts of electricity of positive and negative charge can be created following the snow-surface tribocharging contact event. Contaminants such as particles tend to move in relationship to thermal and electric field gradients—both of which are present in snow cleaning.\nThe backside of the quartz is typically opposite in charge (conservation of charge) during snow cleaning, therefore the particles once lifted from a front surface migrate around the substrate within a thin-film of subcooled atmosphere and become attracted to the oppositely charged surface on the backside. Quartz cannot be grounded and commonly used antistatic chemical agents contained in the cryogenic cleaning agent would leave stains during cleaning.\nA photoelectric effect has been advantageously employed in different arts for decades. In certain commercial ionization applications, the photoelectric effect is used to produce highly energetic photons from 0.13 to 0.41 nm (9.5 to 3 KeV) to ionize an atmosphere surrounding a substrate during a production process.\nAs such there is a present need to provide an alternative dense fluid spray cleaning and separation apparatus and process which overcomes the limitations of conventional dense fluid spray technology and provides an environmentally-safe cleaning and finishing alternatives to organic solvents.\nAs such there is a present need to provide clean and effective electrostatic control method during sublimable cleaning processes. Since electrostatic charging is most prevalent in cryogenic cleaning such as carbon dioxide, argon or liquid nitrogen blasting—a three-dimensional ionization method and device is needed to resolve electrostatic charging effects in complex substrates being cleaned, regardless of composition, shape and size."}
{"text": "1. Field of the Invention\nThis invention relates generally to earth-embedded heat exchangers and particularly to such heat exchangers in which an aqueous primary heat transfer fluid is in thermal contact with an encapsulated material with a high latent heat of transition which is capable of undergoing a phase change under the normal operating conditions for a water-to-air heat pump system.\n2. Description of the Prior Art\nWith the increased use of heat pumps in the heating and cooling of buildings has come the recognition that the efficiency of an air-to-air heat pump is high only when the outside ambient temperature is moderate. At ambient temperature extremes, the coefficient of performance of a heat pump falls drastically. At an ambient temperature of 32.degree. F., operation continues only at an energy loss because the evaporator must be defrosted. As a consequence, alternate heat sources which would remain essentially constant despite fluctuating ambient air temperatures have been sought.\nGround water, under favorable conditions, is one such source; and ground water-to-air heat pumps have been operated at a high coefficient of performance year-around. However, the water sources commonly used in the past have been either well-water or city water. Either source would be quickly exhausted if it were widely adopted as a heat source/sink. These drawbacks can be overcome with the use of a closed water loop in which the earth itself serves as a giant heat sink for air conditioning in the summer and as a heat source from which large amounts of heat warmer than the ambient air can be extracted in the winter. Unfortunately, attempts to utilize heat pump systems in which the heat exchange medium is circulated in a closed loop through the earth have never worked very well.\nProblems with such closed loop systems arose because little of the available heat in the earth's crust can be abstracted with the heat exchangers of the prior art. These heat exchangers typically comprised long sections of copper coil. No enclosures surrounding these coils and containing phase change materials were provided. Due to the low thermal conductivity and heat capacity of the earth, the energy in the vicinity of such earth coils was rapidly dissipated when the instantaneous heating demand of a heat pump system was placed on one of them. As a consequence, the temperature of the heat transfer fluid continued to fall; and within a relatively short period of time, a second earth coil which had been idle for a lengthy time interval and which had had sufficient time to recover had to be substituted for the exhausted earth coil. Alternately, the heat pump system had to be shut down.\nSuch earth coils, in addition to being very expensive and costly to install because of the long sections required to effect adequate heat transfer between the earth and the heat transfer medium, also tended to separate from the surrounding earth upon freezing during winter operation and to overheat plant roots in lawns and gardens as a result of summer heat dissipation. Moreover, a substantially large open area in close proximity to the structure to be conditioned was required in which to embed the earth coils; at least one sq. ft. of open area was required for every sq. ft. of heated floor area.\nA further deterrent to the widespread use of such closed loop systems was the need for custom-designed components due to the variation in soil properties from one locale to another. This obstacle has been compounded by the addition of thermal storage tanks, the most commonly employed means in the prior art for reducing the instantaneous heating or cooling demand of the heat pump system upon the earth coil."}
{"text": "1. Field of the Invention\nThe present invention relates to a snowmobile ski, and particularly, a snowmobile ski having an offset runner and keel and a snowmobile ski having longitudinally extending snow compacting ribs.\n2. Description of Related Art\nIn designing snowmobile skis, there are several performance criteria, which are deemed to be important. These include, inter alia, maneuverability, steerability, stability, flotation, and side-hilling capability. xe2x80x9cManeuverabilityxe2x80x9d is also referred to as xe2x80x9caggressivityxe2x80x9d in the industry and describes the ability of the skis to provide an effective and superior handling/steering of the sled in response to the turning of the steering bar of the snowmobile. xe2x80x9cSteerabilityxe2x80x9d refers to the relative ease with which the sled can be steered with the skis. By definition, a set of skis is easy to steer when the resistance to the steering is kept minimal. xe2x80x9cStabilityxe2x80x9d addresses the pitching and darting tendency of the snowmobile skis during the ride. Rather than being aligned in the intended direction without much wavering, snowmobile skis may pitch and dart, thereby requiring frequent adjustments of the steering system in order to steer the skis and the sled in the intended driving direction. xe2x80x9cFloatationxe2x80x9d refers to the performance of the skis on powder snow (as opposed to compact trails or ice) which allows the sled to be propelled over snow. In such applications as mountain snowmobiles, superior floatation characteristics are essential, as the snowmobile will often be driven on deep fresh powder snow. Finally, xe2x80x9cside-hilling capabilityxe2x80x9d refers to the performance of the snowmobile skis, typically on mountain snowmobiles, when the sled climbs a hill by traversing diagonally up the hill in a zigzag fashion. In side-hilling, the snowmobile is tilted in the direction of the turn, and as a result, the ski on the tilted side tends to be relied upon more heavily than the ski on the other side during each turn. In instances where the skis have poor side-hilling capabilities, the snowmobile may slip laterally because not enough snow is engaged by the skis when tilted on one side.\nDepending on the intended application of the snowmobiles, the particularities of the designs of the snowmobile skis or the methods of attaching the skis to the snowmobile can be varied so as to optimize the desired performance characteristics of the skis. Indeed, prior patents illustrate numerous examples of attempts to improve desirable characteristics of the snowmobile skis.\nFor example, U.S. Pat. No. 5,038,882 (the xe2x80x9c\"\"882 patentxe2x80x9d) shows a conventional snowmobile ski which comprises a ski liner adopted to be connected to a bridge on its top surface, a keel disposed longitudinally along the bottom snow-contacting surface of the ski liner, and a runner (also referred to as a xe2x80x9cwear rodxe2x80x9d in the art) disposed on the keel. The bridge connected to the ski liner\"\"s top surface is adopted to connect to the ski leg of the snowmobile. The ski shown in the \"\"882 patent is typical of conventional skis in that the keel extends longitudinally along the middle of the ski bottom surface. The \"\"882 patent states that a more comfortable ride can be obtained by a spindle arrangement which offsets its longitudinal axis from the ball joints connected to each spindle. Such offset spindle arrangement, according to the \"\"822 patent, induces the ski to xe2x80x9ctoe-outxe2x80x9d thereby providing improved stability, i.e., less pitching and darting.\nIn other examples, U.S. Pat. Nos. 5,040,818 and 5,145,201 disclose a snowmobile ski having a plurality of concave surfaces extending along the snow contacting bottom surface of the ski. The ski shown in these two patents has a center concave surface extending along the middle of the ski and two concave steering surfaces, which is recessed from the center concave surface. The \"\"818 and \"\"201 patents state that the plurality of concave portions provides floatation by compacting snow and lifting the ski and that the lateral recessed steering portions provide better stability and control as the ski turns. However, the center concave surface extends on the bottom surface of the keel, which longitudinally extends along the middle of the ski, as done in other conventional skis. Similarly, while U.S. Pat. No. 3,643,979 shows a snowmobile ski having recessed lateral surfaces or stepped portions, the ski in the \"\"979 patent also has a runner and a keel longitudinally extending along the middle of the ski in a conventional manner.\nWhile snowmobile skis disclosed in the aforementioned patents have keels and runners longitudinally extending along the middle of the skis, several other prior patents disclose snowmobile skis with the keels or runners at the lateral sides of the skis, thereby forming a shape of a tunnel or channel. For example, U.S. Pat. Nos. 3,734,221 and 3,817,544 disclose a snowmobile ski having a longitudinally extending central groove that creates a concave shape. The runners are provided on both sides of the groove, and the central groove is said to provide improved steerability. Finally, U.S. Pat. No. 5,836,594 discloses a double-keel snowmobile ski having the keels extending down from the two lateral edges of a flat middle portion of the ski liner, wherein vertically and horizontally extending wedges are provided at the longitudinal ends of the keel and a concave front portion of the ski liner is provided to funnel the snow to the middle portion of the ski liner. According to the \"\"594 patent, a snowmobile ski with the double-keel arrangement disclosed therein enhances floatation, efficient turning, and smooth forward and reverse movement.\nThese and other prior art snowmobile skis heretofore are available, however, do not teach or suggest a snowmobile ski having an offset runner and keel with longitudinally extending snow compacting ribs disposed slanting upward from the keel toward the inner lateral edge of the ski. Snowmobiles skis having these features in accordance with the present invention provide improved performance characteristics beneficial for use in mountain snowmobiles as well as touring type snowmobiles.\nIt is an object of the present invention to provide a novel snowmobile ski with improved performance characteristics for use in both powder snow applications (which is generally beneficial for mountain snowmobiles) and trail applications (which is generally beneficial for touring snowmobiles). It is a further object of the present invention to provide a snowmobile ski with excellent maneuverability, stability, steerability, floatation and side-hilling characteristics.\nIt is an aspect of the present invention to have a ski body with an inner lateral edge, an outer lateral edge, a top surface, a bridge integrally molded with the top surface, and a bottom surface adapted for moving over snow. The inner lateral edge is defined as being the edge, which faces towards the centerline of the snowmobile and the outer lateral edge is defined as being the edge which faces away from the centerline of the snowmobile.\nAccording to an aspect of the present invention, a snowmobile ski has a bridge adapted to connect to a ski leg of a snowmobile, and the ski body with the inner lateral edge and the outer lateral edge. The top surface of the ski body is adapted to be connected the bridge, and the bottom surface of the ski body is adapted for moving over snow.\nA single keel is provided lengthwise along the bottom surface of the ski body in such a manner that a centerline of the keel is offset from a centerline of the bridge in a direction toward the outer lateral edge of the ski body. Finally, a runner is connected to the keel on one axial side of the ski body.\nAnother aspect of the present invention provides a snowmobile ski with a longitudinally extending rib arranged on the snow contacting bottom surface of the ski body between the inner lateral edge of the ski body and the keel, which is disposed lengthwise along the bottom surface of the ski body in such a manner that a longitudinal axis of the keel is offset from a longitudinal axis of the bridge in a direction toward the outer lateral edge of the ski body.\nIn accordance with yet another aspect of the present invention, a snowmobile ski has a plurality of longitudinally extending snow compacting channels arranged on the bottom surface of the ski body. In the preferred embodiment, each of these snow compacting portions extend from the ski body to different levels, but as illustrated in a second embodiment, the compacting portions could be situated at the same level. When viewed cross-sectionally along the transverse direction of the ski, the preferred arrangement of the snow compacting portions are arranged such that the bottom surface of the ski is generally inclined in a direction from the keel toward the inner lateral edge.\nIt is still yet another aspect of the present invention that a snowmobile ski has a ski body with an inner lateral edge and an outer lateral edge where the snowmobile ski has a keel which extends downwardly from the bottom surface of the ski body at the outer lateral edge and only one runner which is attached to the keel.\nAccording to another aspect of the invention, a snowmobile ski has a ski body with an integrally molded bridge. The integrally molded bridge is adapted to be connected to a ski leg of a snowmobile. Also molded into the bridge are indentations to accept steel bushings, which are press-fitted into the indentations to ensure a solid surface between the ski body and the ski-leg. The bushing are designed such that head of the nut and bolt holding the ski-leg to the bridge do not increase the width of the ski. A single keel is provided lengthwise along the bottom surface of the ski body in such a manner that a centerline of the keel is offset from a centerline of the bridge and or the centerline of the ski width, in a direction toward the outer lateral edge of the ski body. The centerline of the ski is determined by a line, which separates the ski body in a left side and a right of the same width. Finally, a runner is connected to the keel on one axial side of the ski body.\nIn another aspect of the invention, along each side of the longitudinal rib disposed on the bottom surface of the ski body is defined snow compacting portions which extend from the ski body at different levels, such that when viewed cross sectionally, the bottom surface of the ski generally inclines in a direction from the keel toward the inner lateral edge.\nA flat section between the keel and the outside lateral edge of the ski improves maneuverability and reduces snow projection.\nOther objects and advantages of the present invention will be apparent from the appended claims and understood from the following detailed description of the present invention and the following drawings."}
{"text": "Electroluminescent lamps typically provide approximately 30 foot-lamberts of illumination and are thus suitable for various low-intensity illumination applications, such as decorative lighting, egress lighting, cockpit and dashboard display panels, and membrane switches. They have also been used as backlighting sources for liquid crystal display (LCD) devices. However, most LCD applications, including black/white and color LCD displays and high definition displays, require greater backlighting illumination than electroluminescent lamps can provide. Furthermore, most electroluminescent lamps have poor maintenance characteristics: they typically degrade to about half their initial brightness within 2000 hours of operation.\nFluorescent lamps, which provide between 2000 and 4000 foot-lamberts of illumination, have been used as illumination sources for these LCD display devices. However, when used in LCD display applications, fluorescent lamps have their own disadvantages. For example, they are bulky. Being made of glass, they are also fragile and thus are unable to withstand rugged environments; if broken, they may release small amounts of mercury. They also do not operate at temperatures below -20.degree. C. In contrast, electroluminescent lamps do not have the disadvantages of size and construction that fluorescent lamps have. They are quite small and thin, light in weight, extremely rugged, and they can operate at temperatures well below -20.degree. C.\nTo be useful in LCD backlighting applications, electroluminescent phosphors must emit in narrow bands of the blue, green and red zones of the visible spectrum. Specifically, the blue emission wavelength should preferably be between 460-470 nm, the green emission wavelength should preferably be between 535-545 nm, and the red emission wavelength should preferably be between 610-650 nm.\nOne electroluminescent phosphor which meets the above-described blue emission color requirements for LCD backlighting applications is a blue-emitting copper-activated zinc sulfide phosphor, ZnS:Cu. Zinc sulfide phosphors and methods of making them are described in U.S. Pat. Nos. 2,807,587 to Butler et al., 3,031,415 to Morrison et al., 3,031,416 to Morrison et al., 3,152,995 to Strock, 3,154,712 to Payne, 3,222,214 to Lagos et al., 3,657,142 to Poss, and 4,859,361 to Reilly et al., all of which are assigned to the assignee of the instant invention. However, notwithstanding their emission color characteristics, none of these electroluminescent zinc sulfide phosphors, nor any other electroluminescent phosphors, are sufficiently bright for use in most LCD backlighting applications or high definition display devices.\nIt would be an advantage in the art to provide an electroluminescent phosphor having improved luminance and maintenance for use in LCD and high definition display devices, and a method of making the phosphor."}
{"text": "Deduplication is a technique commonly used in the computer storage industry to reduce the storage footprint by removing common portions of files. Deduplication is typically performed by scanning entire files, creating chunks of smaller size within the scanned files, and thereafter creating a fingerprint for the byte stream in each of the chunks within a file. The fingerprint for each chuck of a file is then compared against a fingerprint database. The fingerprint database stores fingerprints for previously processed chunks of files.\nIf a particular fingerprint created for a chunk (“chunk X”) is not present in the fingerprint database, then that particular fingerprint is unique (i.e., it has not been previously encountered) and the fingerprint, and information about its association to chunk X, is stored in the fingerprint database. On the other hand, if the fingerprint is already present in the fingerprint database, then this means that chunk X is considered to be a duplicate of another chunk (“chunk Y”) already represented by another fingerprint stored in the fingerprint database; consequently, a pointer reference is saved that associates the previously stored fingerprint (which is also associated with chunk Y) to chunk X."}
{"text": "1. Field of the Invention\nThe present invention relates to a substrate cleaning apparatus and a substrate cleaning method that clean a substrate, in particular, a semiconductor wafer.\n2. Background of the Invention\nTypically, a resist on a semiconductor wafer substrate is removed (cleaned off) by immersing the semiconductor wafer in a chemical (resist remover) in a remover tank of a substrate cleaning apparatus.\nIn a conventional cleaning technique, for example, an SPM (sulfuric acid-hydrogen peroxide mixture) resist removing process, a mixture (SPM cleaning fluid) containing sulfuric acid and hydrogen peroxide in a certain proportion is first produced. Typically, to improve the resist removing capability, the SPM cleaning fluid serving as a remover is kept at a high temperature equal to or higher than 100 degrees Celsius, for example. Generally, hydrogen peroxide (H2O2) thermally decomposes into water (H2O) and oxygen (O2). In addition, hydrogen peroxide has a boiling point of 150 degrees Celsius and evaporates when heated to such a high temperature. Hydrogen peroxide tends to thermally decompose and evaporate at high temperatures in this way, so that the oxidation capability (resist removing capability) of the SPM cleaning fluid is lowered on the time passes.\nAccording to an existing measure against this, before introduction of a semiconductor wafer or the concentration of the SPM cleaning fluid is managed, a predetermined amount of hydrogen peroxide solution is additionally injected into an inner tank (wafer process tank) or an outer tank (overflow receiving tank) of a process tank at desired intervals, thereby maintaining the oxidation capability of the SPM cleaning fluid.\nA method of managing the hydrogen peroxide concentration uses a concentration monitor. However, since the concentration has to be always kept constant, the amount of usage of the hydrogen peroxide solution and sulfuric acid increases substantially.\nAccording to the method in which the hydrogen peroxide solution is introduced before introduction of the semiconductor wafer described above, it takes long before the hydrogen peroxide solution is sufficiently mixed with the circulating mixture to form a uniform SPM mixture containing the hydrogen peroxide solution and sulfuric acid.\nIn the case where the hydrogen peroxide solution is introduced into the inner tank as described above, the hydrogen peroxide solution having a low specific gravity cannot be effectively mixed with sulfuric acid and overflows into the outer tank, and thus, a desired resist removing capability cannot be obtained.\nSimilarly, in the case where the hydrogen peroxide solution is introduced into the outer tank, since the hydrogen peroxide solution is not effectively mixed, the hydrogen peroxide solution having a low specific gravity is not mixed with the sulfuric acid and remains separated in the upper part of the outer tank. Thus, it takes a predetermined time before a uniform mixture of sulfuric acid and the hydrogen peroxide solution is formed, or it is difficult to feed all the introduced hydrogen peroxide solution to a circulation piping because some of the hydrogen peroxide solution is discharged from the overflow pipe rather than being fed to the circulation piping.\nFurthermore, for example, a conventional substrate cleaning apparatus using the SPM cleaning fluid has a process tank that holds a mixture containing a hydrogen peroxide solution and sulfuric acid and is used for cleaning a substrate immersed in the mixture, circulation piping that extends between a primary side of the process tank on which the mixture is injected into the process tank and a secondary side of the process tank on which the mixture is discharged from the process tank and has a pump for causing circulation of the mixture, a filter disposed in the circulation piping for removing particles in the mixture, a chemical tank for additionally injecting a mixture of a hydrogen peroxide solution and sulfuric acid into the process tank, and a heater for heating the mixture in the chemical tank (see Japanese Patent Laid-Open No. 6-342780, for example).\nHowever, in the conventional substrate cleaning apparatus, the mixture is heated by the heater to a temperature higher than a temperature required in the process tank before the mixture is injected into the process tank. Thus, the SPM cleaning fluid heated to the high temperature thermally decomposes and evaporates, so that the hydrogen peroxide concentration of the SPM cleaning fluid decreases. As a result, a desired SPM cleaning fluid cannot be obtained in the process tank.\nAs described above, the conventional technique described above has a problem that the SPM cleaning fluid cannot have a required resist removing capability, and the amount of the hydrogen peroxide solution additionally injected into the SPM cleaning fluid cannot be optimized."}
{"text": "Various electronic devices such as a smart phone, a tablet PC, a portable multimedia player (PMP), a personal digital assistant (PDA), a laptop personal computer (PC), and a wearable device may provide various functions (for example, social network service (SNS), internet, multimedia, photograph and performance of photo moving images) as well as a telephone function. In particular, various electronic devices provide a function of reproducing content (for example, media moving image).\nUsers watch content using an electronic device under various living environments. Accordingly, the electronic device provides a function of automatically or manually controlling brightness of a display depending on the surrounding environment. However, the existing electronic devices may uniformly control the brightness of the display, or the like independent of content."}
{"text": "The present invention generally relates to image forming systems, such as printers, copiers and facsimiles, and in particular to both side image forming systems capable of printing images on both the sides of paper or the like.\nIn a printer employing electrophotography printing, such as a laser printer, video data is formed per page on a memory region by using source data (font data of a character, image data and so on) and page data (data formed per page on the basis of data supplied from a host computer system). Then the video data thus formed is read out from the memory and a corresponding image is printed on a sheet of paper.\nCurrently, a printer having a both side printing function is available. In a both side printing process which achieves the both side printing function, a front side of paper is subjected to the printing operation and the paper is passed through a path for turning the paper over. Then, a rear side of the paper is subjected to the printing of the next page. In the conventional printer, the path for turning the paper over (paper inversion path) is considerably long. Therefore, it becomes possible to print another paper while the preceding paper is being turned over in the paper inversion path. With this printing process, an increased printing speed can be obtained.\nHowever, in order to attain the above both side printing operation, it is necessary to intermittently skip a page to be printed. It is therefore required to form video data amounting to a few of pages in the memory in advance. For this requirement, a memory having an extremely large capacity must be provided in the printer. In particular, a printer capable of handling a sheet of paper of a large size such as A3 becomes expensive.\nThe above fact holds true for facsimiles and copiers employing the electronic photograph printing method."}
{"text": "The present invention relates to a system for transmitting a digital video signal and recording the received video signal. More particularly, the present invention relates to great extension of the range of use of a digital signal recording/reproducing system by greatly shortening a recording time through transmission of a video signal in a compressed form, and further relates to great extension of the range of use of a digital signal recording/reproducing system by making the number of signals to be recorded and a recording/reproducing time variable.\nAs a digital magnetic recording/reproducing system (hereinafter referred to as VTR) is conventionally known, for example, a D2 format VTR. In such a conventional digital VTR, the elongation or shortening of a reproducing time is possible by using variable-speed reproduction. However, the prior art reference does not at all disclose high-speed recording in which a recording time is shortened to 1/m, multiple recording in which a plurality of signals are recorded, and the compression/expansion of a recording/reproducing time.\nThe above-mentioned conventional digital VTR has a feature that a high quality is attained and there is no deterioration caused by dubbing. However, the shortening of a dubbing time is not taken into consideration. Therefore, for example, in the case where a two-hour program is to be recorded, two hours are required. Thus, there is a drawback that inconveniences are encountered in use. Also, the multiplexing of recording signals is not taken into consideration. Therefore, for example, when two kinds of programs are to be simultaneously recorded or reproduced, two VTR's are required. This also causes inconveniences in use."}
{"text": "Nuclear receptors constitute a large superfamily of ligand-dependent and sequence-specific transcription factors. Members of this family influence transcription either directly, through specific binding to the promoters of target genes (see Evans, in Science 240:889–895 (1988)), or indirectly, via protein-protein interactions with other transcription factors (see, for example, Jonat et al., in Cell 62:1189–1204 (1990), Schuele et al., in Cell 62:1217–1226 (1990), and Yang-Yen et al., in Cell 62:1205–1215 (1990)). The nuclear receptor superfamily (also known in the art as the “steroid/thyroid hormone receptor superfamily”) includes receptors for a variety of hydrophobic ligands, including cortisol, aldosterone, estrogen, progesterone, testosterone, vitamin D3, thyroid hormone and retinoic acid, as well as a number of receptor-like molecules, termed “orphan receptors” for which the ligands remain unknown (see Evans, 1988, supra). These receptors all share a common structure indicative of divergence from an ancestral archetype.\nLipophilic hormones such as steroids, retinoic acid, thyroid hormone, and vitamin D3 control broad aspects of animal growth, development, and adult organ physiology. The effects of these hormones are mediated by a large superfamily of intracellular receptors that function as ligand-dependent and sequence-specific transcription factors. The non-steroidal nuclear receptors for thyroid hormone (TR), vitamin D3 (VDR), all-trans retinoic acid (RAR), and fatty acids and eicosanoids (PPAR) form heterodimers with the 9-cis retinoic acid receptor (RXR) that bind bipartite hormone-response elements (HREs) composed of directly repeated half sites related to the sequence AGGTCA (Mangelsdorf and Evans, Cell 83: 841–850, 1995). In contrast, the steroid receptors function as homodimers and bind to palindromic target sequences spaced by three nucleotides (Beato et al., Cell 83: 851–857, 1995). In addition to the known receptors, a large group of structurally-related “orphan” nuclear receptors has been described which possess obvious DNA and ligand binding domains, but lack identified ligands (Mangelsdorf et al., Cell 83:835–839, 1995; Enmark and Gustafsson, Mol. Endocrinol. 10:1293, 1996); and O'Malley and Conneely, Mol. Endocrinol. 6:1359, 1992)). Each has the potential to regulate a distinct endocrine signaling pathway.\nIt is widely viewed that the hormone response is a consequence of the release, from an endocrine gland, of a ligand that circulates through the blood, and coordinately regulates responses in target tissues by acting through specific nuclear receptors. Hormone responsiveness is dependent on the ability to rapidly clear ligand from the blood and the body so that, in absence of a stimulus, target tissues return to a ground state. Hormonal homeostasis is thus achieved by the coordinated release and degradation of bioactive hormones. Steroid hormones and their many metabolites are primarily inactivated by reduction and oxidation in the liver. Since hundreds of adrenal steroids have been identified (e.g., dozens of each of the sex steroids (androgens, estrogens and progestins), 25–35 vitamin D metabolites, and likely hundreds of fatty acids, eicosanoids, hydroxyfats and related bioactive lipids), the problem of efficient ligand elimination is critical to physiologic homeostasis. In addition to the existence of a myriad of endogenous hormones, a similar diversity of ingested plant and animal steroids and bioactive xenobiotic compounds must also be degraded. Such compounds often are lipophilic and may accumulate to toxic levels unless they are metabolized to water-soluble products that can be readily excreted. Therefore, the efficient detoxification of harmful xenobiotics is essential to the survival of all organisms.\nSelye first introduced the concept that exogenous steroids and pharmacologic substances may function to modulate the expression of enzymes that would protect against subsequent exposure to toxic xenobiotic substances (H. Selye, J. Pharm. Sci. 60:1–28, 1971). These compounds, which Selye called “catatoxic steroids,” are typified by the synthetic glucocorticoid antagonist, pregnenolone-16-carbonitrile (PCN). PCN, and a variety of xenobiotic steroids, induce the proliferation of hepatic endoplasmic reticulum and the expression of cytochrome P450 genes (Burger et al., Proc. Natl. Acad. Sci. (USA) 89:2145–2149, 1992; Gonzalez et al., Mol. Cell. Biol. 6:2969–2976, 1986; and Schuetz and Guzelian, J. Biol. Chem. 259:2007–2012, 1984). Cytochrome P450 (CYP) enzyme(s), present in the endoplasmic reticulum of livers, often catalyze the initial step in such detoxification pathways. P450's are crucial for the detoxification of most xenobiotics, including various environmental pollutants, procarcinogens, and drugs (for review see Denison M S and Whitlock Jr, J. Biol. Chem. 270:18175–18178, 1995). In addition, CYPs are also responsible for the reduction and oxidation of steroid hormones and their many metabolites.\nOne consequence of PCN treatment is the induction of nonspecific “protection” against subsequent exposure to such diverse xenobiotics as digitoxin, indomethacin, barbiturates, and steroids (Selye, supra, 1971). Furthermore, it is known that a variety of such compounds can activate P450 genes responsible for their detoxification or degradation (Fernandez-Salguero and Gonzalez, Pharmacogenetics 5:S123–128, 1995; Denison and Whitlock, supra 1995; O. Hankinson, Ann. Rev. Pharmacol. Toxicol. 35:307–340, 1995; and Rendic and Di Carlo, Drug Metab. Rev. 29:413–580, 1997). P450's constitute a superfamily; each form possesses an overlapping but distinct substrate specificity. Some P450 genes are expressed constitutively, while others, particularly those involved in xenobiotic metabolism, are inducible. In many cases, inducers are also substrates for the induced enzymes, therefore, P450 activities typically remain elevated only as needed. Among the CYP gene family members, the CYP3A isoenzyme is of particular significance from a medical perspective. The human CYP3A4 enzyme is involved in the metabolism of a large number of clinical drugs including antibiotics, antimycotics, glucocorticoids, and the statin class of HMG-CoA reductase inhibitor (Maurel P, Ioannides C Ed. (CRC Press, Boca Raton, Fla.,). pp. 241–270, 1996). Indeed, the drug-induced CYP3A4 activation constitutes the molecular basis for a number of important clinically known drug drug interactions. CYP3A23 and CYP3A11 are rodent homologues of CYP3A4 in rat and mouse, respectively. Indeed, purified CYP3A11 (P450MDX-B) exhibited comparable activity to CYP3A1 (another rat CYP3A homologue, Halvorson, et al., Arch. Biochem. Biophys. 277:166–180, 1990) and CYP3A4 (Yamazaki, and Shimada, Arch. Biochem. Biophys., 346:161–169, 1996) for testosterone 6β-hydroxylation, which is thought to be one of specific reactions for the CYP3A enzyme in rodents and primates (Matsunaga et al., 1998). The regions of the 5′ regulatory sequences of CYP3A23 and CYP3A11 share high homology, including multiple putative response elements (Toide et al., Arc. Biochem. Biophy. 338:43–49, 1997), indicating similar transcriptional regulatory mechanisms among these rodent CYP3A genes.\nAlthough there are substantial structural and catalytic similarities among the various members of the CYP3A family across species lines, important differences exist in regulatory control of these genes (for review, see Gonzalez F J., Pharmacol. Ther. 45:1–138, 1990., and Nelson D R., Arch. Biochem. Biophy., 369:1–10, 1999). For example, a clear discrepancy between human and rodents is that the antibiotic RIF induces CYP3A4 in human liver (Watkins et al., N Engl J Med 338:916–917, 1985) but does not induce CYP3A23 in rats (Wrighton et al, Mol Parmacol 28:312–321, 1985) and CYP3A11 in mice (Schuetz et al., Proc Natl Acad Sci USA 93:4001–4005, 1996), respectively. On the other hand, the anti-glucocorticoid PCN, which induces CYP3A23 in rat liver (Wrighton et al, 1985), only weakly induces human CYP3A4 (Schuetz et al., Hepatology 18:1254–1262, 1993, Kocarek et al., Drug Metab Dispos 23:415–421, 1995, Blumberg et al, Genes Dev 12:3149–3155, 1998), and does not induce CYP3A6 (Dalet et al., DNA 7: 39–46, 1988), a rabbit homolog with a drug response specificity similar to CYP3A4 (Barwick et al., Mol Pharmacol 50: 10–16, 1996). Given the widespread metabolic importance of CYP3A, it would be of great clinical benefit to find an appropriate animal model for use in developing a better understanding of the regulatory control and inter-individual heterogeneity in liver expression of CYP3A in humans.\nWhile it appears that catatoxic compounds such as PCN regulate the expression of cytochrome P450s and other detoxifying enzymes, two lines of evidence argue that such regulation is independent of the classical steroid receptors. First, many of the most potent compounds (e.g., PCN, spironolactone, and cyproterone acetate) have been shown to be steroid receptor antagonists; whereas others (e.g., dexamethasone) are steroid receptor agonists (Burger, supra, 1992). Second, the nonspecific protective response remains after bilateral adrenalectomy (and presumably in the absence of adrenal steroids), but not after partial hepatectomy (Selye, supra, 1971).\nInsight into the mechanism by which PCN exerts its catatoxic effects is provided by the demonstration that PCN induces the expression of CYP3A1 and CYP3A2, two closely related members of the P450 family of monooxygenases (see, for example, Elshourbagy and Guzelian in J. Biol. Chem. 255:1279 (1980); Heuman et al., in Mol. Pharmacol. 21:753 (1982); Hardwick et al., in J. Biol. Chem. 258:10182 (1983); Scheutz and Guzelian in J. Biol. Chem. 259:2007 (1984); Scheutz et al., in J. Biol. Chem. 259:1999 (1984); and Gonzalez et al., in J. Biol. Chem. 260:7435 (1985)). The CYP3A hemoproteins display broad substrate specificity, hydroxylating a variety of xenobiotics (e.g., cyclosporin, warfarin and erythromycin), as well as endogenous steroids (e.g., cortisol, progesterone, testosterone and DHEA-sulfate. See, for example, Nebert and Gonzalez in Ann. Rev. Biochem. 56:945 (1987) and Juchau in Life Sci. 47:2385 (1990)). A PCN response element (which is highly conserved in the CYP3A2 gene promoter) has since been identified in subsequent studies with the cloned CYP3A1 gene promoter (see Miyata et al., in Archives Biochem. Biophysics 318:71 (1995) and Quattrochi et al., in J. Biol. Chem. 270:28917 (1995)). This response element comprises a direct repeat of two copies of the nuclear receptor half-site consensus sequence AGTTCA.\nIn addition to inducing CYP3A gene expression, PCN has also been shown to have marked effects on hepatic cholesterol homeostasis. These effects include significant decreases in the levels of HMG-CoA reductase and cholesterol 7a-hydroxylase gene expression, with associated reductions in sterol biosynthesis and bile acid secretion. PCN has also been reported to enhance the formation of cholesterol esters and the hypersecretion of cholesterol into the bile. Thus, PCN affects key aspects of cholesterol metabolism, including its biosynthesis, storage and secretion.\nActivation of orphan nuclear receptor(s) by catatoxic steroids provides a possible mechanism for the induction of xenobiotic metabolizing enzymes by compounds that do not activate known steroid receptors. Because such enzymes are activated by high (pharmacological) doses of xenobiotic and natural steroids, such a “sensor” would be expected to be a broad-specificity, low-affinity receptor. Such receptors could be activated not only by endogenous steroids and metabolites but also by exogenous compounds such as phytosteroids, xenobiotics and pharmacologic inducers. Indeed, it is known that a variety of such compounds can activate P450 genes responsible for their detoxification or degradation (see, for example, Fernandez-Salguero and Gonzalez in Pharmacogenetics 5:S123 (1995); Denison and Whitlock, Jr., supra, 1995); Hankinson in Ann. Rev. Pharmacol. Toxicol. 35:307 (1995); and Rendic and Di Carlo in Drug Metab. Rev. 29:413 (1997)).\nIn healthy individuals, steroid levels are tightly regulated, with increased catabolism of endogenous steroids being compensated by the pituitary releasing an increase of ACTH, which stimulates biosynthesis, and maintenance of plasma steroid levels. The increased catabolism is reflected by elevated urinary levels of steroid metabolites. Indeed, it is already known that treatment with rifampicin increases urinary metabolites, such as 6β-hydroxycortisol (Ohnhaus et al., Eur. J. Clin. Pharmacol. 36:39–46, 1989; and Watkins et al., J. Clin. Invest., 83:688–697, 1989), and bile acid metabolites, such as 6β-hydroxy hyocholic and 6α-hyodeoxycholic acids (Wietholtz et al., J. Hepatol, 24:713–718, 1996), while the plasma levels of many circulating steroids rise slightly due to increased synthesis (Lonning et al., J. Steroid Biochem. 33:631–635, 1989; Bammel et al., Eur. J. Clin. Pharmacol, 42:641–644, 1992; and Edwards et al., Lancet 2:548–551, 1974).\nWhen synthetic steroids, such as prednisolone (McAllister et al., Br. Med. J. 286:923–925, 1983; and Lee et al., Eur. J. Clin. Pharmaco,. 45:287–289, 1993) or 17α-ethynylestradiol (F. P. Guengerich, Life Sci., 47:1981–1988, 1990) are administered together with rifampicin, plasma levels are rapidly decreased due to enhanced urinary clearance. In some patients undergoing rifampicin therapy for tuberculosis, the increase in urinary steroid levels has led to misdiagnosis of Cushing's syndrome (Kyriazopoulou and Vagenakis, J. Clin. Endocrinol. Metab., 75:315–317, 1992; Zawawi et al., Ir. J. Med. Sci., 165:300–302, 1996; and Terzolo et al., Horm. Metab. Res., 27:148–150, 1995). In these patients, steroid production and clearance normalized when rifampicin was withdrawn. In patients with Addison's disease, who mostly lack the ability to synthesize adrenal steroids, rifampicin treatment leads to rapid depletion of endogenous and administered steroids. These documented clinical situations confirm that induction of CYP3A4 causes increased steroid catabolism (Kyriazopoulou et al., J. Clin. Endocrinol. Metab. 59:1204–1206, 1984; and Edwards, supra, 1974). However, the art is silent regarding the mechanism by which steroid metabolism is regulated in the body.\nAlthough therapeutically administered steroids are beneficial in achieving therapeutic goals, such compounds can, in some cases, increase the overall level of steroids and xenobiotics above physiologically compatible levels in the subjects to whom they are administered. In other cases, the increased level of steroids and/or xenobiotics may linger in the body longer than is therapeutically required. In addition, some subjects are treated with combinations of steroids and xenobiotics that may be administered separately to treat different conditions, but which, in combination, have an additive, or even synergistic, effect known as a drug interaction. In such cases, the patient may be unaware when a physiologically incompatible level of steroids and xenobiotics has been reached, or when an otherwise therapeutic amount of a steroid becomes potentially dangerous due to combined effects of separately administered drugs.\nThiazolidinediones (TZDs) are a new class of oral antidiabetic agents, and have been identified to be the synthetic ligands for peroxisome proliferator-activated gamma (PPARγ) (for reviews, see Spiegelman, 1998, and Wilson and Wahli, 1997). Troglitazone is the first TZD introduced for clinical use. Although troglitazone is effective in reducing hyperglycemia, concern has been raised by several reports of severe hepatic dysfunction leading to hepatic failure in patients receiving the drug (Neuschwander-Tetri et al, 1998, Shibuya et al., 1998, and for a review, see Watkins and Whitcomb, 1998). The mechanism of the liver toxicity by TZDs remains largely unknown.\nAccordingly, there is still a need in the art for the identification and characterization of broad specificity, low affinity receptors that participate in the mediation of the physiological effect(s) of steroids and xenobiotics, particularly when combinations of such compounds disrupt homeostasis or cause drug interaction."}
{"text": "(a) Field of the Invention\nThe present invention relates to a road anti-dazzle device, and more particularly to an anti-dazzle device that includes a cylinder made of high-density polyethylene. The cylinder has a plurality of longitudinal adjacent curved faces that have turbulent functions to enhance the cylinder's resistance to wind pressure. The sectional area of the cylinder appears to be the same when viewed from any angle, thereby providing the largest anti-dazzle area to enhance anti-dazzle effects. A plurality of cylinders can be spaced apart on a safety island. A securing disk is used to press the bottom portion of the cylinder to the ground. A bolt is used to firmly position the cylinder on the ground. When the cylinder is hit, it can automatically restore its original shape and will not break to cause secondary harm.\n(b) Description of the Prior Art\nConventional anti-dazzle plates are made of glass fiber or hard plastics (PE). Once they are hit, they will break into pieces, which may cause another road accident. Anti-dazzle plates are installed on safety islands. As their side views reduce their sectional areas, the anti-dazzle effect is affected."}
{"text": "The present disclosure generally relates to an automated sample workcell and a method for operating an automated sample workcell and, in particular, to an automated sample workcell which is operable to process biological samples and a method for operating an automated sample workcell.\nIn analytical laboratories, in particular clinical laboratories, a multitude of analyses on biological samples are executed in order to determine the physiological state of a patient. Current pre-analytical specimen processors on the market are able to prepare a plurality of biological samples such as blood, urine, cerebral-spinal fluid, saliva, etc. Biological fluid samples are typically contained in open or capped sample tubes. Before a chemical or biological analysis can be performed on a biological sample, a plurality of different pre-analytical processing steps may have to be executed on a sample of a patient. Such processing steps may comprise centrifugation, capping, decapping, recapping and/or aliquotation steps. The processing steps may also comprise adding chemicals or buffers to a sample, concentrating a sample, incubating a sample, and the like. A growing number of those ‘pre-analytical’ steps and procedures are executed automatically by automated pre-analytical sample workcells, also known as ‘automated pre-analytical systems’. The kind of analytical test to be executed on a biological sample is typically specified in a test order.\nOne typical problem is with the physical transport and processing of a sample may not always be synchronized with the assignment of a particular analytical test to the biological sample. For example, a blood sample may be drawn from a patient and collected in a sample tube which is then transported manually, or automatically, to an automated sample workcell at a moment in time when it may not clear which kind of biological and/or chemical analysis, also referred to as ‘analytical test,’ shall be executed on that sample. A physician may have to conduct additional examinations of the patient before he/she can decide which kind of analytical test should be executed on that blood sample of the patient. While the physician conducts the additional examinations, the blood sample of the patient may have already arrive at the automated sample workcell, leaving the automated sample workcell with the question what to do with the sample. According to some laboratory settings, a plurality of samples is received together with a pile of paper-based test orders. In those scenarios, the assignment of electronic test orders to samples may be delayed as the data on the paper-based test orders needs to be entered manually into the lab software before an electronic test order can be requested for a particular sample.\nMost state-of-the-art automated sample workcells do not prepare a biological sample for a particular analysis as long as the test order for the sample has not received. Valuable time is lost, as the automated workcell is not able to process the biological sample at all or is merely able to carry samples not having assigned a test order to a buffering station.\nSome state-of-the-art automated sample workcells do not request a test order but rather determine the type of the sample tubes of a biological sample and process the sample exclusively based on its tube type. For example, in one known system, the automated sample workcell determines the sample container type by image analysis for distributing the containers to different areas or racks. A disadvantage of this approach is that information on the tube type alone is often not sufficient to determine all processing steps necessary to prepare a particular biological sample for a particular analytical test. In addition, the systems are inflexible, because in a lab there may exist much more processing workflows and corresponding test orders than tube types.\nGenerally, state-of-the-art automated sample workcells are either completely test order based (and therefore completely dependent on the receipt of a test order) or are solely based on the determination of the tube type. In the latter case, the workcells are inflexible and often not able to sufficiently pre-process a biological sample for a particular analytical test.\nTherefore, there is a need for an improved automated sample workcell which is operable to process one or more biological samples even in case no test order was received for the samples, thereby avoiding delays and speeding up sample processing"}
{"text": "The present invention relates to a method for applying a coating to a part, such as an anti-friction coating to a piston skirt, using a stencil to apply a predetermined thickness of the coating to a pad for subsequent application to the part surface.\nCurrent methods of coating the outer circumference of pistons include conventional pad printing, rotary pad printing, and direct screen printing. The purpose of the coating in the case of automotive skirt coatings is to eliminate scuffing in the skirt area of engine pistons while improving engine performance and reducing hydrocarbon emissions. Additionally, some pistons, such as air conditioner compressor pistons, use the anti-friction coating as a fluid seal. Bearings and shafts can also benefit from a coating as a result of the enhanced anti-friction and reduced wear characteristics imparted to a coated part.\nDirect screen printing is widely used to coat pistons in the 5 micron to 40 micron range of thickness. The drawback in direct screen printing is that the sharp edges of the piston cut through the screen after approximately 5,000 to 10,000 impressions. Direct screen printing of the piston above approximately 45 microns results in an uneven texture produced by the necessarily wide diameter threads of the screen fabric.\nConventional pad printing and rotary pad printing, as disclosed in U.S. Pat. No. 5,266,142, are severely limited with respect to the thickness of coating that can be applied, typically less than 15 microns per pass, or requiring two passes for approximately 25 microns of thickness, due to the use of an etched clichxc3xa9. The etched clichxc3xa9 does not transfer film thicknesses over 15 microns due to the high surface area and cohesion of the coating to the interior surface of the recess in the clichxc3xa9. Additional passes of the pad increase coating time and slow down the coating process.\nIt is desirable in the present invention to provide a method and apparatus for applying a wide range of coating thicknesses from 5 microns to over 100 microns. It is desirable to provide a method and apparatus which can eliminate the problem associated with printing screens wearing out as a result of sharp edges of the part to be coated cutting through the screen. It is desirable in the present invention to provide precise control of the coating thickness being applied to the part. Increased speed of application, particularly of relatively thick coatings is also desirable. It is desirable in the present invention to reduce or eliminate the deformation associated with screens leading to inaccurate or unacceptable coating applications. It is desirable in the present invention to provide relatively thick coats without uneven or textured coating surfaces.\nWhile the present invention will be described and illustrated with respect to a piston, it should be recognized that many parts can be advantageously coated with the apparatus and method according to the present invention including, but not limited to automotive and small engine pistons, two-cycle engine pistons, air compressor pistons, air conditioner compressor pistons, pump pistons, bearings, shafts, inside diameter of bearings, and parts having surfaces to be coated with geometry that is flat, conical, spherical sections, cylindrical sections, and both inside and outside diameters of rings. Generally, for illustration purposes of the present invention, a piston is described as having an aluminum casting with a machined outer diameter, however alternative metals may be used for the part to be coated without departing from the disclosure of the invention.\nAn apparatus according to the present invention applies a coating of material to a predefined area on a surface of a part. The apparatus comprises a movable stencil having a wall or surface with a predefined aperture corresponding to a pattern of the predefined area to be coated, a doctor blade or squeegee, generically referred to here as a wiper, fills the aperture through the stencil with the coating of material to be applied to the part, and a pad receives the coating of material passing through the aperture of the stencil for applying the coating of material to the part in response to relative movement between the part and the pad. A method according to the present invention is for applying a coating of material to a predefined area on a surface of a part and includes the steps of moving a stencil having a wall or surface with a predefined aperture corresponding to a pattern of the predefined area to be coated with respect to a pad, filling the aperture with the coating of material to be applied to the part with a doctor blade or squeegee, generically referred to here as a wiper, and receiving the coating of material passing through the aperture of the stencil on the pad and applying the coating of material to the part in response to relative movement between the part and the pad."}
{"text": "1. Field of the Invention\nThe field of invention relates to baseball glove apparatus, and more particularly pertains to a new and improved baseball glove wherein the same is arranged to minimize discomfort and pain in use of the structure.\n2. Description of the Prior Art\nBaseball gloves of various types have been utilized in the prior art and exemplified by the U.S. Pat. 4,891,845; 4,987,611; 5,031,239; and 3,898,696.\nThe instant invention attempts to overcome deficiencies of the prior art by employing a baseball glove having a front wall, including an aromatic polyamide liner to dissipate impact in use of the structure and in this respect, the present invention substantially fulfills this need."}
{"text": "Endothelin, which has been disclosed by M. Yanagisawa et al., Nature, 332, pp. 411, 1988, is an endothelium-derived vasoconstrictor peptide composed of 21 amino acids and is considered to act on various organs and tissues including blood vessel, trachea and the like through activation of specific receptors on cell membrane. It has been supposed that this peptide can cause contradiction of smooth muscles and that an excessive secretion thereof may lead to various circulatory diseases such as hypertension, coronary ischemia, encephalopathy, nephropathy, circulation failure of various organs, and asthma.\nIt has been reported that TXA.sub.2 -receptor antagonists and inhibitors of TXA.sub.2 -synthetic enzyme and the like can prevent the increase in intracellular calcium-ion level following the excessive secretion of endothelin. However, there have been no reports about specific antagonists against endothelin so far and, therefore, the development of substances capable of inhibiting various actions of endothelin has been demanded. Under the conditions, the present inventors had made intensive study and found that certain triterpene derivatives extracted from Myrica cerifera L. can specifically compete with endothelin for its receptors (PCT/JP/91/01707, W092/12991, U.S. Pat. No. 5,248,807).\nThe present inventors have made continuous efforts with purposes of developing more endothelin antagonists with higher activity and now found that certain novel triterpene derivatives are useful to achieve said purposes.\nThus, the present invention provides a compound of the formula (I) ##STR2## wherein R.sup.1 is hydrogen or a metabolic ester residue; R.sup.2 is hydrogen or --R.sup.3 --R.sup.4 wherein R.sup.3 is --SO.sub.3 --, --CH.sub.2 COO--, --COCOO--, or --COR.sup.5 COO-- (R.sup.5 is lower alkylene or lower alkenylene), R.sup.4 is hydrogen or lower alkyl or a pharmaceutically acceptable salt thereof.\nAlthough all compounds of formula (I) are highly active endothelin antagonists and useful for purposes of the present invention, those wherein R.sup.1 is hydrogen are preferable and those wherein R.sup.1 is hydrogen and R.sup.2 is --COCH.dbd.CHCO.sub.2 H are more preferable.\nFor purposes of the present invention, as disclosed and claimed herein, the following terms are defined below.\nThe term \"metabolic ester-residue\" in definition fox R.sup.1 means independently an ester-residue which decomposes to reproduce biologically active carboxylic acid in a living body. Examples of metabolic ester-residue include (1-acyloxy) lower alkyl such as picvaloyloxymethyl, acetoxymethyl, 1-acetoxyethyl and the like; (1-alkoxycarbonyloxy) lower alkyl such as 1-(ethoxycarbonyloxy)ethyl, 1-(isopropoxycarbonyloxy)ethyl and the like; and (5-methyl-1,3-dioxolene-4-yl)methyl and the like.\nThe term \"lower alkylene\" means C.sub.1 -C.sub.6 alkylene shown by the formula --(CH.sub.2).sub.n --(n=1-6), preferably C.sub.1 -C.sub.4 alkylene such as methylene, ethylene, trimethylene and the like.\nThe term \"lower alkenylene\" means C.sub.1 -C.sub.6 alkenylene such as vinylene, propenylen, butenylene and the like. Preferable groups are those shown by the formula --(CH.dbd.CH).sub.m --(m=1-3).\nThe term \"lower alkyl\" in definition for R.sup.4 means straight or branched C.sub.1 -C.sub.6 alkyl such as methyl, ethyl, propyl, t-Butyl. The compounds (I) of the present invention may form salts with alkali metals (sodium, potassium etc.), alkaline earth metals (calcium, magnesium etc.), ammonium or an organic base (triethylammonium, trimethylammonium etc.).\nThe compounds (I) of the present invention may form salts with inorganic acids (HCl, H.sub.2 SO.sub.4 etc.), organic acids (benzenesulfonic acid, p-toluene sulfonic acid etc.) and the like.\nThe compounds (I) of the present invention having endothelin receptor antagonistic activities are novel. Although the preparation of them can be carried out using methods known in the art, it can be done efficiently according to the following method.\nThus, the compounds (I) can be prepared by reacting a compound of the formula (V): ##STR3## which has been disclosed in WO92/12991 with a corresponding aldehyde under reaction conditions of Horner-Emmons reaction in the presence of an amine such as triethylamine together with cesium carbonate or lithium bromide. The starting compound (V) can be derived from a substance which is obtainable from Myrica cerifera by extraction in a similar manner to that described in WO92/12991. Briefly, twigs of the plant are extracted for several days at room temperature with a polar solvent (e.g., alcohol such as methanol, ethanol, isopropanol, n-butanol, sec-butanol, tert-butanol; acetone; or acetonitrile). The resultant solution is further extracted with a water-immiscible organic solvent (e.g., chlorinated hydrocarbon such as chloroform, dichloromethane; ethyl acetate; or n-butanol) and the extract is chromatographed on silica gel. The isolated substance is then subjected to chemical modification to obtain Myricerone of the formula (IV): ##STR4##\nThe compound (V) is obtained by reacting the compound (IV) with dimethylphosphonoacetic acid.\nThe compound (V) is then reacted with an aldehyde of the formula (III): ##STR5## wherein R.sup.2 is as defined above and R.sup.6 is t-butoxycarbonyl (hereinafter, referred to as Boc) or hydrogen under the reaction condition for Horner-Emmons reaction. The condensation product is deprotected and/or chemically modified to give the compound (I) of the present invention.\nThe Horner-Emmons reaction can be effected in either of following two methods."}
{"text": "This invention relates to tools used to adjust fuel systems on vehicles, and, more particularly, to a tool that provides to the mechanic a direct indication of the degree of adjustment.\nCarburetors, fuel injectors, and other types of fuel system devices (termed generally \"carburetors\" herein) are found in automobiles and trucks to mix fuel with air in a controlled manner. The mixture flows from the carburetor into the cylinders of the engine, where it is burned to release energy and drive the engine. The operation of the engine therefore depends directly upon maintaining the carburetor in proper adjustment, and carburetor adjustment is a key part of major and minor tuneups.\nThere are two principal points of adjustment on carburetors and fuel injectors, the mixture adjustment and the idle speed adjustment, and each has an independent adjustment control in the form of a screw that can be rotated clockwise or counterclockwise. In tuning the engine, the mechanic alternatively adjusts these two controls while observing engine performance, to obtain an optimal performance of the engine.\nIn modern engines, the carburetors are sometimes rather inaccessible to the mechanic, with the result that it may be difficult to adjust the carburetor controls while observing the engine operation. Elongated adjustment tools have been developed to aid in reaching the controls, but the presently available tools are difficult to use. There is therefore a need for an improved carburetor adjustment tool that provides access to the carburetor, yet is easy to use in performing fine adjustments of the carburetor controls. The present invention fulfills this need, and further provides related advantages."}
{"text": "1. Field of the Invention\nThe invention relates to a bearing element with rolling bodies, with an inner raceway, an outer raceway and adjacent rolling bodies arranged rotatably between the inner and the outer raceways.\n2. The Prior Art\nRotating machine parts generally move at a rotational speed with respect to the surroundings at which they are fastened. The connection is a rotatable bearing which, at high rotational speeds, is customarily a known type of rolling bearing. In order to maintain the rotatability over the long term, the bearings have to be lubricated. In most applications, the outer part of the bearing, the outer bearing ring, is fixed in the housing, and the rotating machine part, generally the shaft, rotates together with or without the inner bearing ring.\nThe differences in speed between the outer and inner bearing rings of the bearing result in a different pressure composition in air boundary layers close to the surface.\nAccording to the continuity equation for flowing media (Bernoulli's law), a higher static pressure prevails at the outer bearing ring, which is a standstill or moves slowly, then at the inner bearing ring which rotates rapidly. This has the effect that, in the boundary layer region of the inner bearing ring, there is a suction effect in the bearing interior while the ambient pressure prevails in the boundary protection region of the outer bearing ring.\nBy means of the centrifugal force of the air layer rotating together with the inner bearing ring, reinforced by the rolling bodies which likewise rotate at the same time, a type of radial fan is formed, as a result of which the air is pressed counter to the pressure drop according to Bernoulli to the outer bearing ring and flows off to the outside.\nAt the laterally sealed bearing point, at which the lateral sealing disk bears tightly against the outer bearing ring, the air flows back radially again in the direction of the inner bearing ring such that a flow component forms in the bearing interior, this flow component, seen by itself, rotating in an opposed manner in the intermediate space between inner and outer bearing rings on both sides of the ring of rolling bodies, but rotating parallel to the axis of the shaft.\nThis flow is combined with the tangential component caused by the rotation, such that the starting point is a helical movement of air on both sides of the bearing point. This flow with the pressure differences associated therewith can be used for transporting lubricant in the bearing."}
{"text": "This invention relates to a novel mounting structure for mounting power semiconductor devices such as triacs or silicon controlled rectifiers, and more specifically relates to a novel insulation mounting structure which enables the mounting of the semiconductor device so that it is thermally coupled to an underlying heat sink but is electrically insulated therefrom.\nPower semiconductor devices have numerous applications. By way of example, they can be used in wall box dimmers of the type shown in U.S. Pat. Nos. 3,746,923 and 3,801,874.\nThe mounting of such semiconductor power devices requires a relatively large heat sink to which the device is coupled so that the device can be properly cooled. It is also usually necessary that the semiconductor device be electrically insulated from the electrically conductive heat sink to avoid the danger of shock to a user who might come in contact with the heat sink. The mounting structure should also be of a nature to be easily manufactured and should have sufficient strength to allow handling of the semiconductor device as when attaching leads to the device terminal without causing mechanical damage.\nSome semiconductor devices are manufactured with a semiconductor die contained within the device which is electrically insulated from its conductive outer surface. Such isolated devices can then be mechanically connected to the heat sink by any appropriate fastener, such as a screw or rivet. Thus, the device is inherently electrically insulated from the heat sink. These devices, however, are more costly than devices which do not have internal isolation and generally have poorer heat transfer characteristics than unisolated devices.\nAnother known method for mounting semiconductor devices employs a metalized ceramic disk which is first soldered to the device and then soldered to a subplate of conductive material which is fastened to the heat sink. The metalizing process and soldering process are costly and such devices tend to be sensitive to the heat of soldering. Additional thermal interfaces are also introduced.\nIt is also known to secure semiconductor devices to their heat sink by means of an epoxy cement, using a thin, electrically insulative cement layer. This process requires great care during the manufacturing process and a high degree of cleanliness for proper mechanical adhesion. Moreover, the adhesives are relevantly poor thermal conductors.\nIt is also known to provide mechanical clamps which will press the semiconductor device against a thin insulation layer inserted between the device and the heat sink. These clamps require a large area and are usually metallic so that they too must be electrically isolated from the device or from the heat sink.\nOther connecting systems are known using mica disks with nylon shoulder washers and screws but these are ineffective because the mica is a poor heat conductor and the nylon shoulder washers will flow under pressure, thus reducing clamping force and heat transfer. An arrangement similar to this is disclosed in U.S. Pat. No. 3,801,874 in which a thin insulation sheet is disposed between a TO-220 outline semiconductor device and the heat sink. The opening in the mounting tab of the TO-220 device is aligned with an opening in the insulation sheet and an opening in the heat sink and a single conductive screw in employed to compress the three components together. An electrically insulative shoulder washer is required to ensure electrical insulation between the conductive screw and the conductive tab of the semiconductor package. A deformable insulation collar fills in the remaining air gaps in the mounting structure. Thus, a large number of separate parts is needed."}
{"text": "1. Field of the Invention\nThe invention is related to the field of communications, and in particular, to communication network architectures and services.\n2. Description of the Prior Art\nTime division multiplex (TDM) networks may be flat networks. In a flat network each switch or device is connected to every other switch or device. FIG. 1 is a block diagram of a typical flat TDM network. When a connection needs to be made from one switch to another switch there is always a direct connection between the two switches. For example when switch 102 needs to connect to switch 106, the connection is made using trunk 110. Sometimes the direct trunk between switches becomes congested by heavy traffic. Dynamic call routing (DCR) is one method that helps overcome congestion problems. A system using DCR monitors the conditions along the trunks. When a path between two switches has a problem, for example heavy congestion or a trunk failure, the system will use an alternate trunk for the connection. For example, when trunk 110, between switches 102 and 106, has a problem, the system may establish the connection from switch 102 to switch 108 along trunk 112, and then from switch 108 to switch 106 along trunk 114. In this way the connection between switch 102 and 106 is completed without using the trunk with the problem. HDC systems typically monitor the CPU occupancy of a switch or the trunk or link occupancy to track usage at a switch or through a trunk.\nPacket based networks don't typically have dedicated links between each network device (i.e. the network is not flat). Packet based networks typically create virtual paths or virtual connections between two IP addresses. A virtual path typically can be created between two network devices in the packet based network even when there may not be a direct physical link between the two network devices. Sometimes the link or connection between two IP addresses are called tunnels, they can also be called routes. Packet based networks use a similar technique to route traffic around congestion or failures in the network structure. Packet based networks may track usage along a virtual connection by counting the number of hops packets take to reach the destination, they may track latency, or they may track percent bandwidth used at the router or port level. Many TDM networks are now connected to packet based networks. Unfortunately there is no sharing of the path conditions between a TDM network and a packet based network.\nTherefore there is a need for a hybrid dynamic connection routing system to share routing information between the two types of networks."}
{"text": "Reflective optical components for use in the extreme ultraviolet spectral range (EUV), which encompasses the wavelength range from approximately 10 nm to approximately 50 nm, can be realized with multilayer mirrors containing a generally periodic layer sequence comprising a plurality of layer pairs. A layer pair generally contains two layers composed of different materials, which should have a maximum difference in their optical constants in the wavelength range provided for the use of the component. At least one of the materials, the so-called spacer material, should have a minimum absorption at the wavelength provided. The selection of the materials for the multilayer mirrors is therefore dependent primarily on the wavelength at which the optical component is intended to be used. In the EUV spectral range, therefore, for in each case a specific wavelength range that is usually only a few nanometers wide, there is an optimum material pairing that guarantees a high reflection on account of the optical contrast of the layer materials.\nIn the wavelength range from approximately 12.5 nm to 14 nm, which is of great importance, in particular, for the development of optical systems for applications in EUV lithography, multilayer mirrors comprising the material pairing molybdenum and silicon are preferably used since there is particularly good optical contrast between these materials in the wavelength range stated. With Mo/Si (molybdenum-silicon) multilayer mirrors it is possible to obtain, for example, a reflection of approximately 70% at a wavelength of 13.5 nm.\nFor the operation of optical systems for EUV lithography, laser plasma sources that emit at a wavelength of approximately 13.5 nm are provided, in particular, as radiation sources. Since the reflection of the overall optical system in EUV lithography is comparatively low due to the plurality of mirrors, EUV radiation sources of this type have to be operated with high powers in order to compensate for the reflection losses that arise in the optical system. In the vicinity of such a high-power EUV radiation source, EUV multilayer mirrors can be exposed to high temperatures. This is the case, in particular, for an EUV multilayer mirror which is positioned close to an EUV radiation source for beam shaping purposes, for example, as a so-called collector mirror.\nAt high temperatures the materials molybdenum and silicon tend, however, toward the formation of molybdenum silicide, in particular MoSi2, and toward interdiffusion processes at the interfaces, as is known, for example, from German Patent Publication DE 100 11 547 C2. Therefore, at high application temperatures, there is the risk of degradation of such multilayer mirrors, which significantly reduces the reflection. In addition to reduction of the reflection, the degradation caused by interdiffusion processes and molybdenum silicide formation is also associated with a decrease in the thickness of the layer pairs, which is also referred to as period thickness. This decrease in the period thickness results in a shift in the reflection maximum toward a shorter wavelength. The function of an optical system based on Mo/Si multilayer mirrors can be considerably impaired or even completely destroyed by such degradation processes.\nIn order to increase the thermal stability of Mo/Si multilayer mirrors, it is known from German Patent Publication DE 100 11 547 C2 to insert a respective barrier layer of Mo2C at the interfaces between the molybdenum layers and the silicon layers. Furthermore, DE 100 11 548 C2 describes the use of barrier layers of MoSi2 for increasing the thermal stability.\nFurthermore, it is known from U.S. Pat. No. 6,396,900 to insert barrier layers composed of the material B4C into Mo/Si multilayer mirrors in order to increase the reflection and/or the thermal stability.\nThe use of such known barrier layers makes it possible to produce layer systems having a high reflection which have an improved thermal stability by comparison with pure Mo/Si layer systems.\nIn the case of Mo/Si layer systems with barrier layers, however, the technological requirements when producing the barrier layers are comparatively high since the thickness of the barrier layers is generally less than 0.5 nm. In particular, it is difficult to deposit a layer sequence with such thin barrier layers on curved substrates.\nThis holds true particularly if the angle of incidence of the EUV radiation varies over the surface of the multilayer mirror and, for this reason, the layer sequence has to have a layer thickness gradient in order to meet the Bragg reflection condition at all locations of the mirror surface."}
{"text": "1. Field of the Invention\nThe present invention relates to the auctioning of services, particularly to the auctioning of document services, and more particularly to the auctioning of document services over an information exchange network.\n2. Description of Related Art\nIn a market driven economy, one of the expectations that customers have is to obtain the best pricing. Traditionally, suppliers of goods and services establish pricing based on whatever the market would bear, based on the economics of supply and demand. Competition among suppliers tends to drive pricing down to some extent. Competition may be effectively created “on demand” by means of an auction process, for the supply side, the demand side or both. For the supply side, customers can typically obtain the lowest prices from auctioning to the suppliers the goods and services demanded by the customers. Some economists believe that the true market price is established via auctioning.\nTraditional auctioning processes are quite effective for transacting goods, for which quality can be determined relatively easily, as compared to services. However, auctioning of services might not be effective to ensure best pricing because of the nature of services. While the lowest price might be obtained from auctioning a service, it might not be the best pricing considering the potential range of quality in the delivery of the services.\nTake for example document services, and more specifically printing services. Print jobs are not commodities like office supplies, toys or Treasury bonds. Each print job is unique and each printer is unique. Needless to mention, buyers want the best price taking into consideration of the quality of the job and services of the printers. Traditionally, to obtain the best price on a print job, a buyer needs to shop around for several printers, requiring one-on-one dealing with each prospective printer. Such shopping around for the best price is time-consuming. The buyer's choice of printers is limited to printers that are located within certain geographical boundaries that can be reasonably explored by the buyer within time constraints. The result is an economically inefficient market, in which resource allocation and pricing are sub-optimal.\nU.S. Pat. No. 5,826,244 to Huberman (fully incorporated by reference herein) describes a process that relies on the Internet to broker document services. According to the Huberman process, a buyer who needs a particular document services job done can provide a request for these services to an online broker. Service providers can bid competitively on the request by submitting bids to the broker, who awards the job to the lowest bidder. Alternatively, instead of simply picking a winner, the Huberman process offers the buyers a choice of several possible winners, based on, for example, a number of winning bids. The customer can then be given an opportunity to select from among these candidates, or to decline or cancel. It appears that the customers do not preset at the start of the auction the number of winning bids that would be offered to the customers. While the Huberman process improves the efficiency of document services transaction and pricing practices related thereto, its auctioning process has several drawbacks.\nWhile buyers want to encourage printers to compete on price, buyers generally cannot agree to accept the lowest bid. A buyer needs to select the best printer for a requested job based on factors such as the printer's reputation or geographic proximity to the buyer. Not every job and every buyer involves the same tradeoff between price and quality. In some cases, only price matters, but in other cases, jobs may be relatively insensitive to price. By showing the buyers several bids in accordance with the Huberman process, without some control as to the number of bids, the printers realize that they can bid a higher price and may still be able to win the job. Consequently, the Huberman process does not encourage the printers to give their best price while preserving the buyers' ability to select a printer on factors other than price.\nFurther, buyers do not want to be committed to execute a job at all, while printers want to spend their time bidding on real jobs that actually will happen. The Huberman process allows the buyers to rely on their sole discretion to decline or cancel a print job, but it does not provide sufficient protection for the printers to avoid wasting resources to create an estimate for the print job to bid at the auction which would be eventually declined or cancelled by the buyers. Without some level of comfort for the printers to feel that they are not wasting their time participating in the auction process, some printers may quickly lose interest in the auctioning process. Without sufficient printer participation, auctioning may be less effective in establishing lowest pricing to the buyers. Huberman clearly realized that forcing buyers to execute on a job request would not work, as buyers often need to know the cost before they can decide whether or not to execute. Charging a fee to buyers for posting job requests would eliminate the non-serious buyers, but would also discourage serious buyers from auctioning their jobs. Disclosing the buyer's identity to the printers would allow the printers to avoid the buyers known to “habitually” decline or cancel on the print job, but would also encourage the printers to selectively bid on only the jobs of well-known buyers, thereby negating the intended effect of auctioning. Also, buyers worry that their competitors could guess their business plans if their identities were disclosed. For example, if a buyer auctions a job online to print newspaper inserts, its competitor might guess that it plans a new promotion if its identity were published.\nAccordingly, it is desirable to provide an improved auctioning process that will encourage the service providers to give their best prices while preserving the buyers' ability to select a provider based on factors other than price. Further, it is desirable to provide an improved auctioning process that will encourage the service providers to allocate resources to participate in the auctioning process while protecting buyers' confidentiality."}
{"text": "High pressure boilers of the type used by electrical generating plants operate at a water pressure generally in the range of 2,000-4,000 psi. Before such a high pressure boiler can be fired up, it must be supplied with water under pressure, for example, on the order of 500 to 1,000 psi, depending upon the boiler design. Water under pressure is supplied by a series of feed pumps, one feeding the other. Initially, the boiler and the series of feed pumps are usually filled with condensate from a condenser using only the condensate pump. In normal operation, the condensate pump usually takes water from a condenser and increases the pressure to about 150 psi and supplies a condensate booster pump which boosts the pressure to approximately 300-600 psi. In turn, the condensate booster pump supplies water to a boiler feed pump which increases the pressure to 1,000 to 4,000 psi depending upon boiler design and the operating condition, such as start-up, part load, or full load.\nA very conventional arrangement for boiler feed pumps is to have two boiler feed pumps, one being a start-up pump that is limited in size and driven by a constant speed motor, without a fluid drive, and a second separate main “full-size” pump that is used for normal operation and is driven by a variable speed power source, either (a) a mechanical drive steam turbine, (b) a variable speed fluid drive that is in turn driven by the main turbine-generator, (c) a variable speed fluid drive that is driven by a large constant speed electric motor, or (d) a motor driven by a variable frequency power source based on solid state electronics. When a pump is used for boiler feed pump service and it operates a constant speed, the water flow is controlled by a discharge flow control valve (sometimes called a pressure control valve).\nFor boiler feed pump service, it is common to use a two-pole motor, and for 60 hz systems, such motors rotate generally at 3600 rpm if it is a synchronous motor, or between 3575 to 3585 rpm if it is an induction motor (3000 rpm for 50 hz systems). Another motor design that is also commonly used is a four-pole motor, and for 60 hz systems, such motors rotate at or near 1800 rpm (1500 rpm for 50 hz systems), but these motors typically use a step-up gear to increase the pump speed to the 3600 rpm range, or higher, depending upon the pump design.\nAnother conventional arrangement is to have two main pumps “usually approximately 60% capacity each”, that are each driven by mechanical drive steam turbines, wherein for start-up, steam from another boiler, either a dedicated start-up boiler, or a boiler of another operating unit, is used to provide steam to drive one or both of these mechanical drive steam turbines during the start-up phase of this unit. In some of these plants where there are two main boiler feed pumps each driven by a mechanical drive steam turbine, a smaller boiler feed pump with discharge flow control valve is driven by a constant speed motor for start-up, for a total of three pumps. The advantage of this arrangement is that the boiler and turbine-generator can be started using electric power either from the grid or from a “black-start” generator, so that no steam source is needed. Clearly, there are advantages to being able to start using a motor driven by a “black-start” generator located at the plant."}
{"text": "Virtual machines are increasingly being used by enterprises for access to applications and other functions. With increased usage, however, come management and security problems. Issues with virtual machine sprawl (also known as “image sprawl” or “virtual sprawl”) currently provide challenges to the efficient and secure usage of virtual machines. Examples of this include an inability to keep track of who is using which version of a virtual machine and/or an inability to determine whether the correct virtual machines are being used. For example, it is fairly trivial to create a virtual machine for a particular application; yet, this feature may cause management problems because a variety of virtual machines may be created and actively used. This may lead to problems with compliance and change management. The proliferation of virtual machines may make it difficult to keep track of licensing requirements. This proliferation may also make it hard to keep track of how widespread security vulnerabilities have an effect in a deployment infrastructure. Further, derivative virtual machines may lead to similar problems. In addition, applying updates to a deployment of virtual machines may be problematic due to virtual machine sprawl."}
{"text": "The advantages of combinatorial chemistry for the rapid generation of chemical compounds for pharmaceutical screening are well established. The single greatest strength of combinatorial chemistry is that it makes possible the generation of libraries of huge numbers of compounds in a relatively short period of time.\nSeveral methods have been employed for the preparation of such libraries, e.g., solution-phase synthesis; solid-phase synthesis on polymer beads or other divided supports; synthesis on soluble, precipitatable polymer supports; and synthesis on planar supports. All of these methods are capable of generating large libraries (i.e., libraries containing a large number of compounds), but none of them are amenable to rapid screening of the libraries for binding activity, biological activity, or other desirable properties.\nCurrent methods for screening libraries include those in which individual library members, or small groups of members, are assayed in microtiter plates, e.g., by screening for a desired activity or for binding to a specific binding partner, such as a receptor or antibody or other ligand, but the number of compounds that can be assayed at once is on the order of 102 to 104: one compound per well in a 96-well microtiter plate screens at most 96 compounds, while twenty compounds per well in a 384-well plate screens about 7680 compounds. The latter approach, while permitting higher throughput, requires secondary screening to identify the active species in any given well. For example, screening a library containing all of the 3.2 million possible pentapeptides which could be made from the twenty natural amino acids (i.e., a 3.2-million-member library) by these methods would require 500 to 3400 plates. Screening a library of the 64 million possible hexapeptides would require 10,000 to 68,000 plates. Robotic systems are available for microtiter plate assays, but screening a large library by such a method remains a massive undertaking. Furthermore, in a pharmaceutical discovery environment this represents only one of the many assays an organization might wish to conduct.\nAn alternative to testing individual compounds is the testing of complex mixtures, with various “deconvolution” strategies being employed to deduce the active species. These strategies have in common the re-synthesis and re-testing of successively less-complex mixtures. In addition to the great effort involved, the testing of complex mixtures is limited by the low concentration of any individual species in the mixture, and is susceptible to false positive results from the additive effects of large numbers of weakly active species. In practice, deconvolution has had limited success in probing large libraries of compounds. See L. Wilson-Lingaro, J. Med. Chem., 39, 2720-2726 (1996); and D. A. M. Konings, J. Med. Chem., 39, 2710-2719 (1996), and references therein, for a discussion of deconvolution strategies.\nIdeally, the probing of a combinatorial library would be conducted in a single operation, with the active members of the library being in some way “pointed out” of the vast population of compounds by the assay. Two of the current methods meet this requirement. Both methods employ solid-phase synthesis, and both require that the library members remain attached to the solid support. In one method, a library of compounds bound to polystyrene beads is prepared by the split-mix method. The library is assayed in a single batch, by being exposed to a molecule of interest, such as a receptor, enzyme, or other specific binding partner. Any beads to which the molecule binds are visualized (e.g., by a colorimetric assay), and beads so identified are selected and the structure of the library member attached to the bead is determined. This can be done by sequencing if the compound is a peptide or nucleic acid, as described, e.g., in U.S. Pat. No. 5,382,513 (incorporated herein by reference). The structure of the compound on the bead may in some cases be deduced from spectroscopic evidence (see, e.g., U.S. Pat. No. 5,382,513), or by decoding a chemical tag that reveals the chemical history of the bead, as described in patent application WO 95/24186 (incorporated herein by reference). The method is in principle capable of screening very large libraries, limited only by the number of beads one is willing to examine. In practice, libraries of 104 to 106 members can be dealt with in this fashion.\nThe second approach involves physically locating a compound or compounds in a spatially addressable array of compounds on a planar support. In this approach, a compound's identity is revealed by its location in the array. One method of this type employs an array of compounds generated by light-directed synthesis, as first disclosed by Fodor et al. in Science, 251, 767-773 (1991), in which a fraction of sites on a planar support carrying photo-detachable protecting groups is exposed to light through a photolithographic mask, and the fraction of sites thus deprotected are functionalized with a specific monomer or building block, itself carrying a photo-detachable protecting group. The process is repeated with the mask in a different position or orientation, or with a different mask, and a second monomer or building block is attached to the support and/or to the first monomer residues. After numerous such cycles, with careful attention to the pattern of masking, an array of compounds is built up on the support. The final array is completely deprotected, and exposed to the ligand of interest. Binding of the ligand is visualized by immunofluorescence, using antibodies against the ligand which are tagged with a fluorescent dye. Under a fluorescence microscope, any location in the array to which the ligand has bound is visible as a fluorescent area, and the x-y coordinates of the area reveals the identity of the library member to which the ligand was bound. This technology, as applied to polypeptide and oligonuceotide synthesis, is known as Very Large Scale Immobilized Polymer Synthesis, or VLSIPS. It is described in U.S. Pat. Nos. 5,143,854, 5,413,939, 5,424,186, and 5,527,681, all of which are incorporated herein by reference.\nThe photolithographic method of synthesis, however, is cumbersome, and requires a substantial investment in very specialized equipment. A further investment in a fluorescence microscope or a specialized scanner is required for the assay, and highly skilled technicians are required, at least to conduct the synthesis aspect of the process. The scale of library synthesis is limited by the size of the masks and by the translational reach of the scanning device, which together limit the accessible surface for synthesis to a few square centimeters. In practice this technique is presently limited to arrays of 104 to 105 compounds. For these reasons the method is not routinely employed; see G. Jung and A. G. Beck Sickinger in Angewandte Chemie, 31, 367 (1992).\nAs an alternative to photolithography, the use of directed laser light to conduct light-directed synthesis has been described in U.S. Pat. Nos. 4,719,615 and 5,318,679 (both of which are incorporated herein by reference). In the latter patent, a rectangular array support is either held stationary or translated, while a laser beam is scanned across the array by means of a rotating mirror, in the manner of an ordinary laser printer. This provides an alternative to the photolithographic masking approach, but fluorescence microscopy is still relied upon as an assay method.\nAnother alternative method for the synthesis of spatially addressable arrays utilizes ink-jet printing technology to spray micro-droplets of reagent solutions onto a substrate. This method is disclosed in U.S. Pat. Nos. 5,474,796 and 5,449,754, both of which are incorporated herein by reference. This method is in theory capable of preparing arrays of 107 compounds, but a method of indexing such large arrays is not disclosed. Again, fluorescence microscopy is the preferred means of conducting an assay.\nThere remains a need for reliable methods of generating very large, very high-density arrays of chemical compounds, on the order of about 108 or more compounds, along with a method of rapidly screening such arrays for chemical properties of interest, such as binding to antibodies, cellular receptors or other ligands, catalytic activity, or inhibition of enzymes."}
{"text": "Hitherto, there have been proposed magneto-optical recording/reproducing apparatuses adapted to selectively use, as a recording medium for information signals, the magneto-optical disc capable of carrying out overwrite recording of information signals and the reproduction only (type) optical disc.\nThe magneto-optical disc or the optical disc used for a magneto-optical recording/reproducing apparatus as described above is accommodated within the disc cartridge in order to prevent damage of at least the signal recording surface and/or attachment of dust thereto, and is loaded into the magneto-optical recording/reproducing apparatus in the state where it remains to be accommodated within the disc cartridge.\nThe magneto-optical recording/reproducing apparatus (unit) using the disc cartridge comprises holder for holding the disc cartridge, and is adapted so that this holder is rotated from the rising position for carrying out insertion or ejection of the disc cartridge toward the falling position side for carrying out recording and/or reproduction with respect to the magneto-optical disc within the disc cartridge, whereby the disc cartridge inserted and held within the holder is loaded into the recording/reproduction unit.\nThe recording/reproduction unit for carrying out recording and/or reproduction with respect to the magneto-optical disc within the disc cartridge comprises a magnetic head disposed above the holder and adapted to be rotated following rotation of the holder, and an optical head disposed at the position opposite to this magnetic head. Further, at the recording/reproduction unit, there is provided a cartridge loading section within which disc cartridge held within the holder is loaded after undergone positioning, and there is provided a disc rotational operation mechanism for allowing the magneto-optical disc or the optical disc accommodated within the disc cartridge loaded at the cartridge loading portion to undergo rotational operation.\nIn order to carry out recording of information signals such as music signal, etc. by the magneto-optical recording/reproducing unit constituted as described above, the disc cartridge within which the magneto-optical disc which permits overwrite recording of information signals is accommodated is held into the holder to rotate this holder toward the falling position side where recording and/or reproduction of information signals are carried out with respect to the magneto-optical disc to allow the disc cartridge to undergo loading with respect to the cartridge loading portion. When the disc cartridge is loaded with respect to the cartridge loading portion, the magneto-optical disc accommodated within the disc cartridge is clamped with respect to the disc rotational operation mechanism. At this stage, the disc rotational operation mechanism is driven to allow the magneto-optical disc to undergo rotational operation, and to allow the magnetic head and the optical head to undergo movement operation in the radial direction of the magneto-optical disc in a manner synchronous with each other. Then, light beams emitted from the optical head are irradiated onto the magneto-optical disc caused to experience rotational operation to heat the portion where light beams have been irradiated so that its temperature becomes equal to the Curie temperature or more, and to apply, to the heated portion, external magnetic field which has been magnetically modulated in accordance with information signal to be recorded from the magnetic head. Thus, recording of desired information signal is carried out.\nMeanwhile, the magnetic head disposed above the holder and rotated along with the holder is caused to become close to the magneto-optical disc in a manner to slidably come into contact therewith when recording of information signals is carried out with respect to the magneto-optical disc.\nMoreover, when reproduction of information signals recorded on the magneto-optical disc or the reproduction only type optical disc accommodated within the disc cartridge is carried out, since it is unnecessary to apply external magnetic field onto the magneto-optical disc or the optical disc, the magnetic head is moved to the position spaced from the magneto-optical disc or the optical disc. Also when the recording/reproduction unit is in stopped state where no recording of information signals is carried out, the magnetic head is moved to the position spaced from the magneto-optical disc.\nWhen there is no necessity of using the magnetic head as described above, the magnetic head is moved to the position spaced from the magneto-optical disc, etc. in order to prevent abrasion of the magnetic head, or the magneto-optical disc or the optical disc.\nMeanwhile, since the magneto-optical discs or the optical discs can carry out high density recording of information signals, discs of which diameter is caused to be small are used. For example, magneto-optical discs or optical discs of which diameter is caused to be about 64 mm are used. Magneto-optical recording/reproducing units themselves using such magneto-optical discs or optical discs which have been compact as recording medium are caused to become compact. Thus, magneto-optical recording/reproducing units which have been compact into portable size are proposed.\nIn the magneto-optical recording/reproducing units permitted to undergo carrying (caused to be portable) of this kind, realization of thinner thickness and more compact structure is required in order to improve convenience of portability.\nMeanwhile, in order to realize thinner thickness structure of the magneto-optical recording/reproducing unit, it is conceivable to reduce movement quantity of the magnetic head which comes into contact with the magneto-optical disc caused to undergo rotational operation by the disc rotational operation mechanism or becomes away therefrom.\nHowever, in magneto-optical recording/reproducing units in which the magnetic head is disposed above the holder for holding the disc cartridge, when the movement quantity of the magnetic head is simply reduced, a portion of the magnetic head is placed in the state close to the portion within the holder when the holder is rotated so that it is located (placed) at the position where insertion or ejection of the disc cartridge is carried out. For this reason, there is the possibility that when insertion/withdrawal of the disc cartridge is carried out with respect to the holder, a portion of the disc cartridge may come into contact with the magnetic head to damage the magnetic head.\nIn addition, when the holder and the magnetic head are close to each other, there is also the possibility that in such cases that impact is applied to the magneto-optical recording/reproducing unit, the magnetic head may come into contact with a portion of the disc cartridge so that it is damaged."}
{"text": "1. Field of the Invention\nThis invention relates to elongate conductors containing conductive polymers, especially ignition cables for automobiles.\n2. Introduction to the Invention\nThe term \"conductive polymer\" is used in this specification to denote a composition which comprises a polymeric component and, dispersed or otherwise distributed in the polymeric component, a particulate conductive filler. Conductive polymers are well known. When the polymeric component is crystalline, many conductive polymers exhibit positive temperature coefficient (PTC) behavior in a temperature range which starts somewhat below the crystalline melting point. In general, the less crystalline the polymeric component, the less likely it is that the conductive polymer will exhibit sharp PTC behavior. When carbon black is the conductive filler, the nature of the carbon black is also an important factor in determining how resistivity varies with temperature (resistivities referred to in this specification are volume resistivities). Conductive polymers can be shaped in any appropriate way, but are preferably melt extruded, and can be crosslinked, e.g. by radiation, after they have been shaped. Conductive polymers have been widely used, for example, in antistatic flooring, as shielding in high voltage cables, and in devices in which current passes through the conductive polymer between metal electrodes, particularly self-regulating heaters and circuit protection devices which use PTC conductive polymers. For further information about conductive polymers, reference may be made for example to U.S. Pat. Nos. 4,237,441, 4,304,987, 4,388,607, 4,545,926, 4,560,498, 4,591,700, 4,724,417, 4,774,024, 4,935,156, 5,049,850, 5,250,228, and 5,378,407, the disclosures of which are incorporated herein by reference.\nIgnition cables in automobiles must meet very stringent requirements. Thus, they must have satisfactory resistance and capacitance characteristics over a wide temperature range, and must retain those characteristics over a period of years in a very hostile physical environment. Many attempts have been made to provide such cables--see for example U.S. Pat. Nos. 2,790,053, 3,991,397, 4,330,493, 4,363,019, 4,375,632, 4,677,418, 4,704,596, 4,748,436, 4,780,700, 4,894,490, 4,970,488, 5,034,719, and 5,057,812, the disclosures of which are incorporated herein by reference. However the cable constructions which have proved commercially acceptable are complicated and expensive to make, and often require, in order to provide the desired combination of resistance, capacitance and size, two or more conductive components which must be applied in separate operations."}
{"text": "The present invention generally relates to light emitting elements, image sensors, and light receiving elements. More particularly, the present invention relates to a light emitting element, an image sensor, and a light receiving element which are suited for use in an image scanner of a facsimile machine, a copying machine and the like.\nIn facsimile machines, there is an increasing demand to clearly transmit fine drawings, photographs and the like. A primary factor which determines the picture quality is the resolution of the image scanner. By use of an image scanner having a high resolution, it becomes possible to transmit fine drawings, photographs and the like with a high picture quality.\nNormally, in the case of the facsimile machine, a document is read by a sensor array which includes one-dimensionally arranged photoelectric conversion elements. The resolution of the sensor array increases as each sensor element confronts a smaller region of the document surface. However, as the sensor element confronts a smaller region of the document surface, a light receiving area of the sensor element consequently becomes smaller. As a result, a quantity of light received by the sensor element decreases and the signal-to-noise (S/N) ratio becomes poor.\nIn order to eliminate the above described problems, it is possible to consider reducing the electrical noise, improving the light sensitivity of the sensor itself and the like. However, it is most effective to improve the luminance of the light which illuminates the document surface by a light source.\nA xenon lamp, a light emitting element (LED) array or the like is used as the light source for the sensor of the facsimile machine. However, these light sources have poor light converging characteristics. For example, in the case of an image scanner having a resolution of 16 lines/mm or 32 lines/mm, the region of the document surface which must be illuminated by the light from the light source is in a range of 30 to 50 .mu.m with respect to a scanning direction of the image scanner. However, the xenon lamp usually has a tube diameter of at least 1 mm, and it is difficult to converge the light to a width of 0.5 mm or less even by the use of a lens system. On the other hand, when the light from the LED array is excessively converged, there is a problem in that the light intensity become irregular because the LEDs are arranged discretely in the LED array.\nWhen an electroluminescence (EL) element proposed in a Japanese Laid-Open Patent Application No.57-7087 is used as the light source, the width of the emitted light is determined by the thickness of a transparent substrate. But when the width of the emitted light is to be set to 30 .mu.m, for example, the transparent substrate must have an extremely small thickness of 30 .mu.m. When making an elongated light source, the substrate may warp due to a stress generated between the substrate and the EL element, and it is difficult to hold the substrate during a photolithograhy process and the like because the substrate is extremely thin. In addition, a metal reflection layer is used to trap the light. Unlike a total reflection, the reflectivity of the metal reflection layer is not 100%. But when the reflection is repeated a large number of times, the light quantity decreases and there is a problem in that the luminance of the light becomes small at an emitting edge for the light which is generated at a distant location from the emitting edge.\nOn the other hand, another EL element is proposed in Kun et al., \"TFEL Edge Emitter Array for Optical Image Bar Applications\", SID 86 DIGEST, pp.270-272. However, this EL element has an emitter layer which has a thickness in the order of 1 .mu.m and is too thin for use as the light source for the sensor. If the thickness of the light emitting layer were increased to 10 .mu.m, for example, a driving voltage of the EL element must be increased to a range of several kV to several tens of kV, and the reliability of the EL element becomes poor.\nA Japanese Laid-Open Patent Application No.64-89280 proposes an EL element which is combined with a waveguide. According to this EL element, the light emitting layer is provided above or below the EL element. However, the emitted light cannot reach the substrate edge while satisfying the conditions for total reflection because the light emitting element is not provided within the waveguide, and the light intensity inevitably decreases.\nOn the other hand, a light emitting element having a cladding layer provided on both sides of the light emitting layer is proposed in a Japanese Laid-Open Patent Application No.1-109694. According to this structure, it is possible to prevent the light from leaking outside the light emitting layer. However, the luminance of the light emitted from the edge inevitably decreases due to the reabsorption of light within the light emitting layer.\nAs described above, it is essential that the sensor which reads the document has a high resolution in order to obtain a high picture quality. The resolution of the sensor may be affected by stray light existing between a light illuminating system and a light receiving system of the sensor.\nThe light which illuminates the document surface is originally reflected depending on the tone of the document and the reflected light reaches the light receiving element. However, a portion of the light emitted from the light source is reflected at an interface of a thin film which is interposed between the document and the light receiving element, and this portion of the light reaches the light receiving element directly. This portion of the light is referred to as the stray light.\nFIG. 1 is a cross sectional view for explaining the stray light generated in the conventional sensor. FIG. 1 shows a light source 100, a document 101, a substrate 102, a photoelectric conversion element 103, stray light 104, and interfaces 105 of thin films.\nWhen the stray light exists, the photocurrent of the light receiving element is no longer solely dependent on the tone of the document, and the resolution of the sensor deteriorates.\nFor example, Japanese Laid-Open Patent Applications No.61-100073 and No.58-106947 propose arrangements for preventing deterioration of the resolution due to the stray light, by providing an optical waveguide in the light illuminating system and the light receiving system. However, according to these proposals, there are problems in that the light source and the waveguide must be assembled, the production process is complex because of the need to form the light receiving element on the edge of the waveguide so as to cover the edge and the like. For these reasons, it is difficult to reduce both the size and production cost of the sensor.\nNext, a description will be given of an edge receiving type image sensor (image reading element) which uses an optical waveguide in the light receiving system. Compared to the image sensors of the conventional facsimile machines using a reduction optical system or a full-size imaging system, the image sensor which uses the optical waveguide is advantageous in that a light illuminating system for illuminating the document and the light receiving system for receiving the light at the photoelectric conversion part can be isolated optically. For this reason, the image sensor which uses the optical waveguide is developed for applications where a high resolution is required, such as the case where fine graphics, drawings, photographs and the like are to be read.\nIn the edge receiving type image sensor, the light received from the edge of the optical waveguide reaches the photoelectric conversion part, and an incident angle of the light at the edge of the optical waveguide is uniquely determined depending on a ratio of refractive indexes of a core layer and a cladding layer. In other words, when the ratio of the refractive indexes of the core layer and the cladding layer of the optical waveguide is small, the incident angle of the light which is trapped within the optical waveguide and is transmitted to the photoelectric conversion part is small. On the other hand, when the ratio of the refractive indexes of the core layer and the cladding layer of the optical waveguide is large, the incident angle of the light which is transmitted to the photoelectric conversion part is large. Accordingly, in the edge receiving type image sensor, it cannot be determined unconditionally whether the ratio of the refractive indexes of the core layer and the cladding layer should be large or small. A more detailed description will now be given in conjunction with FIGS. 2A and 2B.\nFIG. 2A shows an example of a case where a ratio of a refractive index n.sub.core of a core layer 211 and a refractive index n.sub.clad of a cladding layer 212 is small. In this case, the incident angle of the light is small as indicated by hatchings, and regions of a document 202 which are read by two mutually adjacent optical waveguide edges of an image sensor 201 do not overlap between the document 202 and the image sensor 201, that is, the edges of the optical waveguides to be more accurate. As a result, the resolution of the image sensor 201 is high, and a change in the resolution of the image sensor 201 with respect to a change in the distance between the document 202 and the sensor 201 is small. But on the other hand, the absolute quantity of the light input to the optical waveguide edges from the document 202 is small because of the small incident angle of the light. For this reason, a large photocurrent cannot be obtained, and the S/N ratio of the image sensor 1 is deteriorated thereby.\nOn the other hand, FIG. 2B shows an example of a case where the ratio of the refractive index n.sub.core of the core layer 211 and the refractive index n.sub.clad of the cladding layer 212 is large. In this case, the incident angle of the light is large as indicated by hatchings, and the absolute quantity of the light input to the optical waveguide edges from the document 202 is large. However, the regions of the document 202 which are read by two mutually adjacent optical waveguide edges of the image sensor 201 overlap between the document 202 and the image sensor 201. As a result, the resolution of the image sensor 201 is poor, and a change in the resolution of the image sensor 201 with respect to a change in the distance between the document 202 and the sensor 201 is large.\nTherefore, according to the conventional image sensor, it is impossible to simultaneously obtain satisfactory resolution and operation stability and satisfactory light quantity which is input to the optical waveguide edges, where the operation stability refers to the stability with respect to the change in the distance between the document and the image sensor."}
{"text": "1. Field of the Invention\nThis invention relates to a method of forming very narrow electrically conductive lines and narrow spacing between such lines on printed wiring boards (PWBs).\n2. Brief Description of the Prior Art\nThe continual decrease in the size of electronic components has resulted in a continuing effort to obtain finer circuit lines and spacing between the lines on printed wiring boards. In the manufacture of certain types of printed wiring board systems or modules, particularly those systems having a pair of printed wiring boards with a core therebetween, sequential copper plating steps are required. Generally, plating first takes place on the initial printed wiring boards with the build-up of copper on the board surfaces and vias or holes through the board electrically interconnecting the two major surfaces. That board and a second board on which a similar plating operation has taken place are then bonded to a core layer positioned between the two boards. The module is plated with copper including vias or holes through the two boards and the core to electrically interconnect the two initial boards. These steps involve the drilling of holes in one or both boards, plating of those holes, bonding of the boards to the core layer to form a module, then drilling additional holes through the entire module with subsequent plating of the additional holes. A result of this sequential plating requirement is that there is a build-up of the plating material on the surfaces of the initial boards from the plural plating steps. This large build-up of copper provides a problem in obtaining fine lines of copper and fine spaces between the copper lines. The final etching of the surface copper is dependent upon the etch factor of the etchant chemistry and the copper thickness on the surface to be etched. A consequence of this problem has been relatively poor yields and a large minimum line and space width.\nAttempts have been made to improve the etch factor through the use of chemical additives or by minimizing the amount of copper plated on the board surfaces. While both of these approaches are valid and have shown improvement in the ability to reduce line widths, they have limited applications. Even with a greatly improved etch factor due to etchant chemistry, the surface copper thickness will dominate line widths. It is therefore apparent that other approaches are necessary in order to provide line widths and spacing between lines of reduced dimensions while improving the yields."}
{"text": "1. Field of the Invention\nThe present invention relates to a handle switch for a vehicle. More particularly to a handle switch for a vehicle including a composite switch which is suitable for operating a navigation display device, etc.\n2. Description of Background Art\nHeretofore, composite switches have been known that are designed for a reduced layout space and improved operability by having a plurality of operating elements centralized in one place.\nJapanese Patent Laid-Open No. 2005-302347 discloses an arrangement wherein a four-way switch is disposed outside of a centrally located push switch and a rotary switch is disposed outside of the four-way switch. Japanese Patent Laid-Open No. 1998-199374 discloses an arrangement wherein a composite switch doubling as a push switch and a four-way switch is disposed centrally and a rotary switch is disposed outside of the composite switch. Japanese Patent Laid-Open No. 2009-4209 discloses an arrangement wherein a composite switch doubling as a four-way switch and a rotary switch is disposed outside of a central push switch.\nIn recent years, there has been a growing demand for vehicles, not only four-wheeled vehicles but also motorcycles, which have a display device or a navigation device that can be connected to the Internet for browsing the web, and which allow the vehicle occupant to operate the display device or the navigation device while in the riding posture when the vehicle is temporarily stopped. The composite switch disclosed in prior art identified above is suitable for operating those devices. However, since the handle of the motorcycle has switches for operating electric devices including a horn, a headlight, etc., the composite switch needs to be reviewed in terms of its structure and layout if it is to be added and also in view of making itself resistant to erroneous operation."}
{"text": "The present invention generally relates a method of controlling the fuel pressure within the fuel rail of an internal combustion engine.\nWithin an internal combustion engine fuel delivery system, a fuel rail delivers fuel to fuel injectors that inject the fuel into the cylinders of the engine. The pressure of the fuel being injected through the fuel injectors is critical, therefore it is important to keep the pressure within the fuel rail as consistent as possible. Closed loop control systems for controlling the fuel pressure within fuel rails have been developed, but do not adapt to system variations such as part-to-part tolerance differences and wear. Therefore, there is a need for an improved method of controlling the fuel pressure within the fuel rail of an internal combustion engine."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of determining the time of modification of the physical state of a fluid medium and, in particular, determining the coagulation time of blood, of the kind in which a ferromagnetic ball is placed in the bottom of a cup containing the fluid to be studied and driven with a periodic movement under the effect of an external magnetic field, and in which the modifications of the movements of this ball due to the modification of the physical state are detected. The invention also relates to a device for implementing this method.\n2. Description of the Prior Art\nSuch a method is described more particularly in the European patent n.degree. 90 192 in which the ball thus driven moves freely over the flat and horizontal bottom of a cup. When the modification of the physical state of the fluid which is contained in the cup, and which is here blood plasma to which reagents have been added, occurs, the ball moves over a curve directed towards the inside because, so it seems, of the increase of mechanical resistance opposed by the contents of the cup, and its radial movement is detected by means well known in themselves.\nThe accuracy of measurements thus made leaves much room for improvement, since the radial movement of a ball towards the center of a cup of very small size is a physical magnitude which is difficult to detect and to evaluate satisfactorily. Furthermore, the reproducibility of such measurements is not good, because, particularly, of the complexity of the phenomena on which they are based."}
{"text": "This invention pertains to insecticides and, more particularly, to an insecticide for controlling whiteflies.\nWhiteflies have caused extensive damage to food and fiber crops in Texas, Florida, California and Arizona. They attack some 500 crop and other plant species as well as poinsettias and other ornamentals. The sweet potato whitefly has been a major pest to many agronomic crops such as cotton, watermelon and cantaloupe. The whitefly also attacks area landscape by overwhelming plantings of lantana and passion vine. Arizona's long growing season along with a extensive list of suitable hoists for sweet potato whitefly make this pest extremely dangerous to agronomic and horticultural plantings. Crops attacked by whiteflies in the USA include: cotton, alfalfa, carrots, peanuts, cantaloupe, citrus, eggplant, honeydew melon, broccoli, cabbage, lettuce, celery, peppers, and vegetables.\nAs if California growers, who are periodically attacked by the Mediterranean fruit fly, didn't have enough to fret about, in the past two years along came another winged peril. A new type of whitefly, Bemisia tabaci, swirled like massive clouds of dandruff across fields of melons, lettuce, broccoli, and other vegetables. The pest was dubbed the poinsettia strain of the sweet potato whitefly in Florida, when it was sighted on ornamental plants. In California, the numbers went from troublesome to devastating in recent growing seasons. Pesticides and other remedies failed. People were even trying to vacuum the insects off their plants. Crop damage totaled 118 million dollars, 2,200 jobs were lost, and a state of emergency was declared in Imperial and Riverside Counties, because of whiteflies.\nPopulations of Bemesia tabaci whiteflies remain relatively light throughout the winter months due to their low fecundity during cooler months of the year but as the temperature rises crops like cantaloupe are planted in the spring and populations of Bemesia whiteflies begin to rise. As the spring melons are harvested the adult whitefly populations begin their migration to cotton where whiteflies mate. Females whiteflies lay their eggs on the underside of cotton leaves. Without grower intervention 6 to 12 generations of whitefly can be produced in a single growing season.\nSuccessful control of sweet potato whitefly depends on destruction of the whitefly hosts. Whitefly control depends on the assistance of farmers and growers in an agronomic environment and the assistance of gardeners and homeowners in towns, cities and suburbs. Since cotton and vegetables are important crops to the agricultural economy, it is unlikely that either crop will be abandoned in an attempt to reduce the host range of the whitefly.\nMany types of insecticides have been used to kill whiteflies and other insect pests with varying degrees of success. Unfortunately, conventional insecticides often pollute the atmosphere, create toxic waste, are not biodegradable, and contaminate rivers, streams, ponds, lakes, the ground, soil, and underground aquifers. Conventional insecticides can also damage crops, gardens, lawns, trees and shrubs. Furthermore, conventional insecticides can be hazardous to food supplies and can be harmful and even toxic to humans, animals, birds and fish.\nIt is, therefore, desirable to develop an improved insecticide which overcomes most, if not all, of the preceding problems."}
{"text": "The present invention relates to connectors containing fusible materials to assist in forming a connection and more particularly to such connectors which, during the heating of the fusible material, form part of a circuit, the temperature of which is autoregulated at about the Curie temperature of magnetic material included in the circuit at least during the heating operations.\nU.S. Pat. No. 3,243,211 discloses a connector containing a fusible material so that upon insertion of an object to be joined to the connector or insertion into the connector of two members to be joined, and upon heating of the connector, the fusible material is caused to melt and contact said object or objects and upon cooling, effect a bond. The connector may also include a heat recoverable member whereby the liquified fusible material is bonded and caused to contact the object or objects while in the fluid state. This device requires an external heat source to melt the fusible material, such as hot air or an infrared radiant source.\nThe difficulty with the device of the patent is in the danger of overheating the objects to be soldered or otherwise bonded as well as adjacent objects. In the electronics art, for instance, overheating of delicate integrated circuits is a problem as is overheating of circuit boards, mastics, resins, heat shrinkable polymers, glues, potting compounds, all of which can be destroyed by the application of excessive heat. Further, the device has little utility for joining wires, tubes or members which are large effective heat sinks, since the large amount of heat required cannot be readily transferred through the heat shrinkable sleeve without damaging it."}
{"text": "Processes to heat treat food products are known in the art. U.S. Pat. No. 4,956,532, for instance, describes a method and apparatus for even and rapid heating, pasteurization or sterilization of products contained in a package, such as, for example, pharmaceutical products or food, which are conveyed through a microwave treatment chamber on a continuous conveyor belt. In order to heat components of the products having different heat absorption for the purpose of reaching the pasteurization temperature rapidly and evenly, the temperature of the components is recorded by a temperature sensor, and on the basis of this the computer calculates the respective ΔT, taking into account given product parameters and the desired temperature to be reached. When the product passes under the following input aperture, a given microwave energy level is coupled into each specific component, as a result of which an even temperature of all products in the package is achieved very rapidly without overheating of the product.\nFurther, WO2011062499 describes for instance a device for pasteurizing a mass of foodstuff which comprises: a feed; a first tube of an electrically and magnetically inert material suitable for contact with foodstuff; an arrangement of electrodes added to the first tube and connected to an RF power generator so that the mass present in the first tube can be heated in this first tube. Co-acting electrodes are disposed with a substantial mutual axial interspace."}
{"text": "Compounds with antimicrobial properties have attracted great interest in recent times as a result of an increase in the prevalence of infections caused by Gram-positive bacteria, resulting in serious or fatal diseases. Furthermore, the regular use of broad spectrum antibiotic formulas has led to the increased occurrence of bacterial strains resistant to some antimicrobial formulations.\nNovel antimicrobial compounds have the potential to be highly effective against these types of treatment-resistant bacteria. The pathogens, having not previously been exposed to the antimicrobial formulation, may have little to no resistance to the treatment."}
{"text": "The present invention relates to an office workwall assembly, and in particular to a workwall assembly that comprises a frame assembly mountable to a wall structure and adapted to support office utilities assemblies thereon, such as shelving assemblies, storage assemblies, and drawer assemblies."}
{"text": "As computer systems have become faster in processing speed, improvement in the reliability of such systems has become an issue.\nThis means that a technique is needed that can make systems operate stably without shutting them down even when a failure has occurred.\nMulti-CPU systems, which are equipped with plural CPUs, can operate stably.\nThe techniques described below are proposed as techniques for the detection of failures in multi-CPU systems.\nPatent Document 1 discloses a technique by which operations are monitored in units of processors. In the technique of Patent Document 1, data is transmitted between processors that are connected via a communications device, and one processor monitors responses from the other processors in order to detect failures.\nPatent Document 2 discloses a technique by which an abnormal operation in a system is monitored in units of instructions so that the system can recover from the failure. Patent Document 2 discloses a method in which the occurrence of a failure in a processor is monitored by the processor itself and the processor performs a resetting operation using an automatic reset occurrence circuit and a failure element holding circuit when the system can recover by resetting.\nThe methods disclosed by the above described patent documents have the problems described below.\nIn the method of Patent Document 1, processors are monitored by communications between processors themselves. This configuration causes a time gap of several minutes between the occurrence of a failure and the recognition of the failure.\nIn the method of Patent Document 1, information used for investigating the cause is collected after the recognition of a failure, and this means that the information collected for the investigation is not exactly from the moment at which the failure occurred. This makes it difficult to use the information for investigating the cause, and thereby the investigation tends to take a long time period or sometimes the investigation itself is prevented.\nThe time scale of monitoring a system (in units of minutes or seconds) and the time scale of a CPU (in units of nanoseconds) are significantly different from each other, resulting in difficulty in understanding the status of the system when a failure occurred, and thereby a long time is required to investigate the cause.\nAlthough there is a technique of collecting information on changes in the status of a CPU in the form of log information, effective information sometimes cannot be collected using this technique because of the limitation on the volume of log information.\nAccording to the technique of Patent Document 2, the monitoring and recovering is performed in units of instructions, and the system is reset in order to restart when the restarting of the firmware can make the system recover from a state involving failure.\nIn the technique of Patent Document 2, the monitoring and recovering is performed in units of instructions, making it difficult to detect a failure in a program or the like in an OS (Operating System)."}
{"text": "The present invention relates generally to an interactive method, system, and apparatus for displaying views and superimposed design style overlays for grooming hair to a desired design style. The invention provides improved visual angels for a surface area plane, and optical sensory digital imaging processing while grooming hair that enables the user to also be instructed on the accuracy of grooming techniques using superimposed overlays interfaced with a camera FIG. 18 (viewing device), having an elongated high impact plastic material embodiment structure, and the hair grooming trimmer 1006 tool as the preferred embodiment in order to achieve a desired hair design style. The invention 102 comprises a system, method, and apparatus for instructional guidance training in order to achieve a desired hair design style using autonomous guidance control system for overlay instructional training.\nAlso, the present invention uses a predictive analytical analysis process of optically determining a change in grooming for comparing hair grooming accuracy based on the superimposed design overlay to guide the user in order to achieve a desired hair design style. According to the present invention, operational control of a grooming hair style design tool using artificial intelligence, and superimposed design styles overlays for grooming hair is embodied, which will allow a sensor to convert image processing to instruct the apparatus controller to operate further allowing the invention to automatically be controlled using microchip 310 embedded processor in order to achieve a desired hair design style. More specifically, this invention relates to the use of various types of grooming tools, such as but not limited to trimmers whereby more accurate grooming of targeted area is achieved.\nWhen a consumer purchases a hair grooming tool such as a trimmer kit; or when visiting a selected hair stylist, the consumer is unaware of the complexities involved with grooming their hair personally; or whether the hair stylist has enough experience to groom the consumer's hair to their desired design style. At home, the consumer must rely\non an inexperienced associate to groom their hair 402 or use a tool such as a mirror FIG. 4 at odd and uncomfortable angles to better gauge the desired method to grooming their hair to a preferred design style as a disadvantage shown in U.S. Pat. No. 7,281,461 and illustrated as prior art in FIG. 17. Referring now to example of prior art being U.S. Pat. No. 7,281,461 FIG. 17, a self-user is shown holding and using the present hair clipper 1702 in various grips and angled positions. In FIG. 17, a self-user is shown holding the hair clipper 1704 in a left hand and using it to clip hair on the left side of the head. In FIG. 17, a self-user is shown holding the hair clipper 1710 in a right hand and using it to clip hair on the right side of the head. In FIG. 17, the blade assembly 1708 carrying the bladeset 1706 is rotated to one of the two positions discussed above as prior art, and in FIG. 17, the blade assembly 1708 carrying the bladeset 1706 is rotated to the other of the two positions discussed above.\nU.S. Pat. No. 5,579,581 assigned on its face to Wahl Clipper Corporation is directed to a clipper blade having multiple cutting edges, namely a cutting edge at each end of the blade. However, the cutting edges on each end are substantially identical such that each blade can be used as either of the fixed blade or the moving blade. Thus, the use of superimposed design overlays for grooming assistance to guide user operations of the same blade assembly is not provided.\nU.S. Pat. No. 5,606,799 also assigned on its face to Wahl Clipper Corporation is directed to a hair clipper having a ball and-socket connection being provided between the handle and the blade assembly. The ball-and-socket configuration allows the blade assembly to be pivoted with respect to the handle. However, the ability to rotate the blade assembly about an axis substantially normal to the cutting plane defined by the blade assembly or a viewing apparatus to view targeted grooming areas is not provided.\nU.S. Pat. No. 5,970,616 is also assigned on its face to Wahl Clipper Corporation. This patent is directed to a hair trimmer that includes a blade housing that is rotatable about an axis substantially parallel to the axis of the handle to vary the angular orientation of the blade housing with respect to the handle. However, the ability to rotate the blade assembly about an axis substantially normal to the cutting plane defined by the blade assembly is not provided. Moreover, the use of intelligent interactive accuracy analysis system is also not provided to afford the user some grooming assistance.\nAlso as discussed above, the blade assembly 1706 is retained in the selected rotational position by the interaction of the lock extension 1708. In this embodiment, a bladeset 1710 is positioned at a particular angle relative to the hair strands to be trimmed. Further, the hair strands are guided toward a cutting zone “Z” of the bladeset 1706 and retained in the cutting zone “Z.”. This, in turn permits a self-user to hold the hair clipper 1704 to position the bladeset at a particular angle in relation to the hair to be trimmed by uncomfortably twisting or pronating and supinating the wrist and forearm, as opposed to bending the wrist sideways at an awkward angle, known as ulnar deviation.\nHowever, if a user performing these repeated un-natural awkwardly angular twists of the wrist without any true; and visually accurate, reference to guide their work, the potential for excessive wear on the wrist could result in decreased stamina or injury in future grooming efforts and assurances that the targeted grooming plane is groomed correctly based on consumer's desired style. When this particular angle of attack of the trimmer or bladeset 1710 relative to the head is substantially a right angle to the hair to be\ntrimmed, a cross-section of the hair presented to the bladeset 1710 to be trimmed is substantially illustrative of the difficulty of grooming hair correctly. Any discussion of the prior art throughout the specification should in no way be considered as an admission that such prior art is widely known or forms part of common general knowledge in the field. Neither solution is desirable to the consumer. Multiple devices are duplicative and expensive while not allowing more intelligent grooming assistance. Thus, it is desirable to provide a grooming device such as a hair clipper or trimmer that permits improved views of targeted grooming area.\nThose skilled in the art should be able to understand the scope of the present invention and the notated prior art patents are representative of this scope as the present invention overcomes the disadvantages:\nTheir manufacture process does not include the use of an imaging device capable of capturing targeted grooming areas on a selected plane, thereby not allowing the consumer to self use the grooming tool in a more efficient and accurate manner while trimming or clipping their or another's hair.\nThe practice of a self use grooming tool being used is hampered by the user not being aware of the instructional steps needed to groom hair to a desired style.\nAccordingly, there is a continuing need for hair clippers capable of providing an intuitive method for grooming hair while using a sensor and imaging processor that present more grooming control to the user. Moreover, there is a continuing need for such a hair clipper or trimming hair grooming device that controls the blade assembly to match operations to the hair design outlined within the superimposed overlays to insure control over quality results when grooming hair.\nThe following is a tabulation of some prior art that presently appears relevant:\nU.S. PatentsPat. No.Kind CodeIssue DatePatentee4,949,460AAug. 21, 1990Sterk4,951,394AAug. 28, 1990MeijerD310,272AAug. 28, 1990Gallanis, et alD310734ASep. 18, 1990Gallanis4,958,432ASep. 25, 1990Marshall5,012,147AApr. 30,1991Bertram, et al5,012,576AMay 07,1991Johannesson5,012,830AMay 07, 1991Vaccaro, et al5,031,315AJul. 16, 1991LabrijnD319324AAug. 20, 1991Wahl, et al5,050,300ASep. 24, 1991Miska5,050,305ASep. 24, 1991Baker, et al5,052,106AOct. 01, 1991Labrijn5,054,199AOct. 08, 1991Ogawa, et alD320673AOct. 08, 1991Sterk5,075,971ADec. 31, 1991McCambridge5,084,967AFeb. 04, 1992Nakagawa, et alD336545AJun. 15, 1993McCambridge, et al5,230,153AJul. 27, 1993Andis5,233,746AAug. 10, 1993Heintke5,235,749AAug. 17, 1993Hildebrand, et al5,245,754ASep. 21, 1993Heintke, et alD339655ASep. 21, 1993SulikD340783AOct. 26, 1993McCambridge5,257,456ANov. 02, 1993Franke, et al5,283,953AFeb 08, 1994Ikuta, et al5,289,636AMar. 01, 1994Eichhorn, et al5,325,590AJul. 05, 1994Andis, et al5,333,382AAug. 02, 1994Buchbinder5,343,621ASep. 06, 1994Hildebrand, et alD351046ASep. 27, 1994McCambridge, et alD352132ANov. 01, 1994Wagenknecht, et alD353022ANov. 29, 1994Simonelli5,353,655AOct. 01, 1994Mishler5,377,414AJan. 03, 1995Buzzi, et al5,383,273AJan. 24, 1995Muller, et alD355506AFeb. 14, 1995Rizzuto, Jr5,398,412AMar. 21, 1995Tanahashi, et al.5,400,508AMar. 28, 1995Deubler5,410,811AMay 02, 1995Wolf, et al.5,414,930AMay 16, 1995Muller, et al.5,423,125AJun. 13, 1995Wetzel5,450,671ASep. 19, 1995Harshman5,463,813ANov. 7, 1995Otsuka, et al.5,469,536ANov. 21, 1995Blank5,473,818ADec. 12, 1995Otsuka, et al.5,479,950AJan. 02, 1996Andrews5,490,328AFeb. 13, 1996Ueda, et al.5,504,997AApr. 09, 1996Lee5,507,095AApr. 16, 1996Wetzel, et al.D368984AApr. 16, 1996Nakashima, et al.D369230AApr. 23, 1996Bone5,524,345AJun. 11, 1996Eichhorn5,530,334AJun. 25, 1996Ramspeck, et al.5,539,984AJul. 30, 1996Ikuta, et al.5,469,536ANov. 21, 1995Blank5,542,179AAug. 6, 1996Beutel5,546,659AAug. 20, 1996Tanahashi, et al.5,548,899AAug. 27, 1996Tanahashi, et al.5,551,154ASep. 03, 1996Tanahashi, et al.5,564,191AOct. 15, 1996Ozawa5,568,040AOct. 22, 1996Krainer, et al.5,568,688AOct. 29, 1996AndrewsD376670ADec. 17, 1996Bellm, et al.5,577,179ANov. 19, 1996Blank5,600,890AFeb. 11, 1997Leitner, et al.5,604,985AFeb. 25, 1997Andis, et al.5,604,986AFeb. 25, 1997Masuda5,611,145AMar. 18, 1997Wetzel, et al.5,611,804AMar. 18, 1997Heintke, et al.D378705AApr. 01, 1997Izumi5,655,301AAug. 12, 1997Dickson5,669,138ASep. 23, 1997Wetzel5,673,711AOct. 07, 1997Andrews5,685,077ANov. 11, 1997Mukai, et al.5,687,306ANov. 11, 1997Blank5,699,616ADec. 23, 1997Ogawa5,701,673ADec. 30, 1997Ullmann, et al.5,704,126AJan. 06, 1998Franke, et al.5,745,995AMay 05, 1998Yamashita, et al5,771,580AJun. 30, 1998Tezuka5,793,188AAug. 11, 1998Cimbal, et al.5,802,932ASep. 08, 1998Vankov, et al.5,842,670ADec. 01, 1998NigoghosianD405923AFeb. 16, 1999Yang5,884,402AMar. 23, 1999Talavera5,884,404AMar. 23, 1999Ohle, et al.D408102AApr. 13, 1999Van AstenD408587AApr. 20, 1999De Visser5,901,446AApr. 20, 1999De Visser5,937,526AAug. 17, 1999Wahl, et al.5,940,980AAug. 24, 1999Lee, et al.5,943,777AAug. 31, 1999Hosokawa, et al.5,956,362ASep. 21, 1999Yokogawa, et al.5,960,515AOct. 05, 1999Lu5,964,034AOct. 12, 1999Sueyoshi, et al.5,964,037AOct. 12, 1999Clark5,970,616AOct. 26, 1999Wahl, et al.5,979,056ANov. 09, 1999Andrews5,980,452ANov. 09, 1999Garenfeld, et al.5,983,499ANov. 16, 1999Andrews6,000,135ADec. 14, 1999Ullmann, et al.6,003,239ADec. 21, 1999Liebenthal, et al.D421818AMar. 21, 2000Mandell, et al.6,044,558AApr. 04, 2000Wu6,052,904AApr. 25, 2000Wetzel, et al.6,052,915AApr. 25, 2000Turner6,079,103AJun. 27, 2000Melton, et al.6,082,004AJul. 04, 2000Hotani6,098,289AAug. 08, 2000Wetzel, et al.D429378AAug. 08, 2000Wu6,125,857AOct. 03, 2000Silber6,126,669AOct. 03, 2000Rijken, et al.6,151,780ANov. 28, 2000Klein6,178,641B1Jan. 30, 2001Meijer6,205,666B1Mar. 27, 2001JunkD439703B1Mar. 27, 2001Wagenknecht, et al.6,219,920B1Apr. 24, 2001Klein6,223,438B1May 01, 2001Parsonage, 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{"text": "In various types of multi-state card games, such as but not limited to blackjack and baccarat, a player competes against an opponent. For purposes of example, the player may be human and the opponent may be a dealer who plays on behalf of a “house” (e.g., a game provider such as a casino). These card games are typically played in hands, where each hand involves an initial number of cards being dealt from a deck to each participant. The cards that each participant holds during a hand may be referred to as the participant's hand, and cards may be added to or taken away from each participant's hand during the course of the hand.\nThe popularity of such card games is based in part on the strategy involved in playing the game, as well as the player being able to risk money or some other form of currency on the outcome of each hand. The possibility of gain and the risk of loss tend to make these games exciting and enjoyable. Nonetheless, for each type of multi-state card game, there are certain combinations of cards that result in a weak hand. When dealt such a hand, more often than not the player will lose to his or her opponent. Being dealt a weak hand reduces the player's level of excitement, as the player's optimal strategy may be to minimize losses rather than to maximize gains. The house that provides a card game usually expects a marginal profit for each hand played. Since exciting card games tend to encourage repeat play, the house profit is likely increased when the overall excitement level of a card game is enhanced."}
{"text": "Identifying moving objects in a scene, from a sequence of images, for example, a video stream, is a very important, and computation-intensive, endeavor. Generally, known approaches to identifying moving objects in pan imagery use “change detection” where an image with moving objects is differenced with, i.e., compared to, a background image that contains no moving objects. The problem is how to generate a background image from a series of images with many moving objects in close proximity to one another.\nThere are several different methods of creating a background image. These methods include: a) for a given frame, taking the previous frame as a background; b) averaging n previous frames to form a background; c) morphing the background by making small changes based on the previous image frame; or d) a 3-frame method using the current frame and the two previous frames.\nA better method of creating a background image that can be used in change detection processing to identify multiple moving objects in close proximity is needed."}
{"text": "ESA antennas are used for a wide range of applications including cellular telephone networks, telemetry systems and automotive, shipboard and airborne radar systems. ESA antennas capable of efficiently radiating over wide bandwidths enable systems having flexibility for multiple mode operation. The growing interest in ultra-wideband (UWB) communications has lead to implementations in which a single ESA antenna is used to accommodate all frequencies of interest. ESA antennas often include an array of notch antenna elements. Each element includes an electrically conductive body having a slot. Generally, the slot includes a feed end which is positioned near a stripline feed and a radiating end which couples the RF signal in the stripline into the air or other medium. The stripline is typically embedded below the surface of a dielectric substrate and extends below the feed end of the slot to enable efficient coupling of an RF signal to be transmitted from the element. The notch antenna element can also be used to couple electromagnetic energy incident at the wide end of the slot into the stripline as a received RF signal. Various parameters affect the frequency content of the RF signal propagating from the element including, for example, the geometries of the base of the notch antenna element and the aperture in a conductive coating on the adjacent surface of the dielectric substrate, and material properties of the dielectric substrate.\nArray antennas constructed of slot antennas and TEM horns generally use vertical feeds that are easily accommodated by a brick architecture as is known in the art. A description of brick architectures and tile architectures is provided in section II of the publication of Robert J. Mailloux, Proceedings of the IEEE, Vol. 80, No. 1, January 1992. Typically, array antennas constructed according to the brick architecture are deeper and heavier than array antennas employing the tile architecture where the distribution of RF signals is accomplished in one or more layers that are parallel to the antenna aperture plane. Conventional notch antennas require a feed that extends away from the antenna element so that layered connections are not practical."}
{"text": "Conventional golf balls can be divided into several general types or groups: (1) one piece balls; (2) two piece balls; (3) wound balls; and (4) other balls with three or more layers. The difference in play characteristics resulting from these different types of constructions can be quite significant.\nBalls having a two piece construction are generally most popular with the average recreational golfer because they provide maximum distance. Two piece balls commonly include a single solid core, usually formed of a crosslinked rubber. Solid cores are often made of polybutadiene that is chemically crosslinked with zinc diacrylate and/or similar crosslinking agents and is covered by a tough, cut-resistant blended cover, such as SURLYN.RTM., an ionomer resin produced by E.I. Du Pont de Nemours & Co. of Wilmington, Del. The combination of the core and cover materials imparts a relatively high initial velocity to the ball which results in improved distance. Due to the rigidity of these materials, two piece balls have a hard \"feel\" when struck with a club.\nAt the present time, the wound ball remain, the preferred ball of the more advanced players due to its spin and feel characteristics. Wound balls typically have either a solid rubber or fluid filled center around which many yards of a stretched elastic thread or yarn are wound. The wound core is then covered with a durable cover material, such as SURLYN.RTM., or a softer material, such as balata or a castable polyurethane. Wound balls are generally softer and provide more spin than the aforementioned two piece balls. Particularly with approach shots into the green, the high spin rate of soft, wound balls enables the golfer to stop the ball very near its landing position.\nRelatively recently, a number of golf ball manufacturers have introduced golf balls having three or more layers in an effort to overcome the undesirable aspects of conventional two-piece balls, such as their hard feel, while retaining their positive attributes, such as increased initial velocity and distance. These balls have multiple core layers, i.e., they include a center with one or more intermediate (\"mantle\") layers, and one or more cover layers. Examples of such multilayer balls include the EPISODE.TM. and the HP2 DISTANCE.TM.(TITLEIST.RTM.).\nAlthough a variety of factors affect which of these types of balls a player will use, all players desire a ball that is affordable and durable. Therefore, in an effort to meet the demands of the marketplace, manufacturers strive to develop low-cost, efficient manufacturing techniques that produce golf balls which are resistant to cutting and cracking, yet which exhibit desirable distance, spin rate, and compression.\nThe durability of a ball depends not only upon its cover, but upon its core as well. A number of elastomers such as polybutadiene, natural rubber, styrene butadiene rubber, and isoprene rubber are commonly used in fabricating golf ball cores. Polybutadiene is most commonly used to obtain desired golfball properties. Manufacturers have added cross-linking agents, such as metallic salts of an .alpha.,.beta.-unsaturated carboxylic acid, to the elastomeric core composition to achieve a desired resiliency, compression, and durability.\nSome manufacturers have instead attempted to provided improved golf balls by surrounding polybutadiene solid centers with thermoplastic mantle layers. For example, U.S. Pat. No. 4,337,946 discloses a golf ball having an intermediate layer of thermoplastic resin between a polybutadiene thread-wound center portion and an outer polyester elastomer cover layer which contributes to the ball's impact and cutting resistance characteristics.\nU.S. Pat. No. 4,919,434 discloses a golf ball having a solid core of more than 40% cis-1,4-polybutadiene and a cover having an inner layer of 0.1 to 2 mm thickness and an outer layer of 0.1 to 1.5 mm thickness. The inner layer is a thermoplastic resin, such as an ionomer, polyester elastomer, polyamide elastomer, thermoplastic urethane elastomer, propylene-butadiene copolymer, 1,2-polybutadiene, polybutene-1, and styrene-butadiene block copolymer, either individually or in combination.\nU.S. Pat. No. 5,439,227 discloses a three-part golf ball having a rubber inner core, and an outer core formed by injection molding a mixture of 100 to 50 weight percent of a polyether ester type thermoplastic elastomer and 0 to 50 weight percent of an ethylene-(meth)acrylate copolymer ionomer.\nWhile materials incorporating a thermoplastic into a polymer blend are known, the use of such blends in portions of a golf ball core is not known. For example, U.S. Pat. No. 4,972,020 discloses an inner cover layer having a modified block copolymer of a thermoplastic polymer and a modified block copolymer consisting essentially of a base block copolymer of a monovinyl substituted aromatic hydrocarbon polymer block and an olefinic compound polymer block having an ethylenic unsaturation degree not exceeding 20 percent, wherein the base block has a molecular unit having a carboxylic acid group and/or a group derived therefrom grafted thereto.\nU.S. Pat. No. 5,093,423 discloses a method of making a thermoplastic elastomer produced by dynamic vulcanization of styrene-butadiene-styrene (\"SBR\") rubber as a dispersed phase of crosslinked SBR, and a co-continuous, matrix of styrene-ethylene-butylene-styrene (\"SEBS\") and polypropylene.\nU.S. Pat. No. 5,100,947 discloses a dynamically vulcanized composition of a polyolefin thermoplastic resin and an elastomer of a rubber material in which a major portion of fillers or specified additives are present in the resin.\nU.S. Pat. No. 5,270,386 discloses a cover blend of vinyl aromatic copolymer and a poly(phenylene ether) concentrate containing poly(phenylene ether), a vinyl aromatic copolymer, polyamide, polycarbonate, polyester, poly(alkyl acrylate), and/or poly(alkyl methacrylate). The blend may optionally contain impact modifiers, thermoplastic molding materials including polyester, polystyrene, polyolefin, polyamide, polyvinyl chloride, polyurethane, polyacetal, and conventional additives, such as dyes and pigments.\nWhile some of the references discussed herein describe the use of thermoplastics in forming a golf ball cover, a golf ball core or portion of a core that contains a blend of both thermoplastic and elastomeric materials is not disclosed. There has thus been a long-felt need, which is now satisfied by the present invention, for a golf ball core, or portion thereof, having a blend of at least one high Vicat softening thermoplastic and at least one elastomer to provide an increased geometric stability without substantially affecting the desired golf ball properties."}
{"text": "Electrical fuses and electrical antifuses are used in the semiconductor industry to implement array redundancy, field programmable arrays, analog component trimming circuits, and chip identification circuits. Once programmed, the programmed state of an electrical fuse or an electrical antifuse does not revert to the original state on its own, that is, the programmed state of the fuse is not reversible. For this reason, electrical fuses and electrical antifuses are called One-Time-Programmable (OTP) memory elements.\nProgramming or lack of programming constitutes one bit of stored information in fuses or antifuses. The difference between fuses and antifuses is the way the resistance of the memory element is changed during the programming process. Semiconductor fuses have a low initial resistance state that may be changed to a higher resistance state through programming, i.e., through electrical bias conditions applied to the fuse. In contrast, semiconductor antifuses have a high initial resistance state that may be changed to a low resistance state through programming.\nContinuous advances in the semiconductor technology oftentimes require changes in the material employed in semiconductor structures. Of particular relevance is the advent of a metal gate electrode, which, in addition to the gate dielectric, a polysilicon layer, and a metal silicide layer, contains a metal gate layer in a gate stack. Typically, the metal gate layer is employed in conjunction with a high-k gate dielectric material. This is because high gate leakage current of nitrided silicon dioxide and depletion effect of polysilicon gate electrodes limits the performance of conventional silicon oxide based gate electrodes. High performance devices for an equivalent oxide thickness (EOT) less than 1 nm require a high-k gate dielectric material and a metal gate electrode to limit the gate leakage current and provide high on-currents.\nThe high-k gate dielectric materials refer to dielectric metal oxides or dielectric metal silicates having a dielectric constant that is greater than the dielectric constant of silicon oxide of 3.9 and capable of withstanding relatively high temperatures, e.g., above 600° C., and preferably above 800° C. The metal gate layer may comprise a metal, a metal alloy, or a metal nitride, and typically has an even higher conductivity than the metal silicide.\nIn view of the above, there exists a need for an electrical antifuse structure compatible with fabrication of other semiconductor devices employing metal gate electrodes and methods of manufacturing the same."}
{"text": "Diamines represent one of the most powerful and versatile class of monomers used in the commercial polymer business. Diamines are widely used in preparing polyimides, polyamides, and virtually all epoxy resins. The use of aromatic diamines is fundamental for preparation of high performance and temperature resistant polymeric materials. For polyimides perhaps one of the best known materials is Kapton® and this material has found widespread use in space applications. For polyamides, Kevlar® represents an aromatic polyamide that has found widespread use in fire-retardant clothing, armor, and many other applications. Thus, aromatic diamines are of great importance for creating thermally-stable/high-performance polymeric materials.\nAlthough aromatic polymers can exhibit resistance to thermo-oxidative processes, it has been demonstrated that in the presence of atomic oxygen (AO) at high kinetic energies that virtually any and all carbon-carbon bonds are destroyed. Thus, Kapton® has been found to rapidly decompose when exposed to AO, both in a simulated environment and in actual low Earth orbit (LEO) testing. While thin coatings of metal oxides and other inorganic materials provide protection from AO erosion, they suffer from defects that occurring during their application along with cracking and failure due to a mismatch in coefficients of thermal expansion (CTE), impact, minor abrasions, or other typical wear and use scenarios. Once a crack appears, undercutting and acute failure of the coating begins leading to catastrophic polymer damage and then material failure.\nOligomeric silsesquioxanes (OS) can be incorporated into a polymer matrix to improve the resistance of the polymer to degradation by AO. In particular, polyhedral oligomeric silsesquioxanes (POSS™) added within a polymer matrix have been shown to greatly enhance resistance of the polymer to attack by AO. There are examples of OS moieties being added/incorporated into high molecular weight and thermally stable polymer materials; however, there are no previous reports where erosion by exposure to AO was completely stopped as in the current invention.\nWright et al. first disclosed the synthesis of a POSS-diamine and its use in preparing polyimide oligomers in the article titled “The Synthesis and Thermal Curing of Aryl-Ethynyl Terminated coPOSS Imide Oligomers: New Inorganic/Organic Hybrid Resins” (Chemistry of Materials 2003, 15(1), pp 264-268). In latter work, Tomzcak et al. used that same POSS diamine to prepare related polyimides that were tested in LEO for durability against AO. Synthesis of the diamine-POSS structure is a multistep process and requires expensive reagents. Furthermore, at high POSS loadings phase separation occurred in the polyimide and unidentified “solids” appeared in the polymer films.\nSvejda et al. in U.S. Pat. No. 6,767,930 describes general methods for incorporating POSS structures into polymer matrices using both reactive methods and guest-host approaches. In this patent the authors demonstrate the value of OS in slowing polymer degradation when exposed to AO.\nLichtenhan et al. in U.S. Pat. No. 6,933,345 put forth several scenarios for POSS incorporation into polymeric materials, both as guest-host and as part of a polymer structure. Specific examples and experimental conditions are not given as well as polymer properties (expected or measured). No methods are given for synthesizing an aromatic OS-diamine or POSS-diamine.\nIn work by Leu et al., “Synthesis and Dielectric Properties of Polyimide-Tethered Polyhedral Oligomeric Silsesquioxanes (POSS) Nanocomposites via POSS-diamine” (Macromolecules 2003, 36, pp 9122-9127) they describe making POSS-polyimides and at relatively low levels of incorporation phase separation (self-assembly) occurred in the polyimide leading to a significant drop in glass transition, increase in coefficient of expansion, and a significant decrease in the material strength.\nLeu et al. also describe an approach in “Polyimide-Side-Chain Tethered Polyhedral Oligomeric Silsesquioxane Nanocomposites for Low Dielectric Film Applications” Chemistry of Materials 2003, 15, pp 3721-3727) by which the POSS structure is incorporated by modifying a fluorinated-polyimide backbone. This requires the polyimide be soluble in organic solvents so lower molecular weight (Mn of 15,000) materials must be utilized. This work was also reported in US patent application 2006/0122350 by Wei et al.\nWright et al. documented new methodology for tethering octahedral silsesquioxanes to soluble polyimides in the article “Chemical Modification of Fluorinated Polyimides: New Thermally Curing Hybrid Polymers with POSS” (Macromolecules 2006, 39, pp 4710-4718). In this paper the authors demonstrate very high loadings of POSS to low molecular weight polyimides and demonstrate that in lab testing films show resistance to erosion by AO. The work requires soluble polyimides, multiple separations of polymer/reactants, and utilizes a specialty monomer that is not commercially available.\nPoe and Farmer in U.S. Pat. No. 7,619,042 (2009) describe a new method for modifying a soluble polyimide material of low molecular weight similar to that described by Wright et al. In the patent they provide a single example of a fluorinated low molecular weight polymer with no characterization of the OS-containing polymeric materials with regard to purity, physical, mechanical, and/or AO resistance.\nIt is to be understood that the foregoing is exemplary and explanatory only and are not to be viewed as being restrictive of the invention, as claimed. Further advantages of this invention will be apparent after a review of the following detailed description of the disclosed embodiments, which are illustrated schematically in the accompanying drawings and in the appended claims."}
{"text": "The storage capacity of hard disk drives has been steadily increasing, as well as the recording density of magnetic recording media mounted therein. In response to these increases, the flying height of a magnetic recording head is becoming smaller over the years, raising the need to reduce the distance between the head and the magnetic recording layer (magnetic spacing) more and more. As one of the methods for reducing the magnetic space, there is a method for thinning a protective layer formed on the magnetic recording layer.\nBecause the protective layer literally serves to protect the magnetic recording layer, the protective layer needs to have a high corrosion resistance in order to prevent the magnetic recording layer, made of metal or the like, from becoming corroded by reaction with atmospheric moisture or corrosive gas. Furthermore, the head moves on the magnetic recording medium at relatively high speed and flies thereon keeping small magnetic spacing therebetween. The protective layer is also required to be durable to the flying travel motion of the head in order to protect the magnetic recording layer from damage even in a case where the head comes into contact with the magnetic recording medium due to some problem. Currently the protective layers are made of DLC (diamond-like carbon) and as thin as 3.0 nm.\nA heat-assisted recording system has been designed as a new recording system for coping with higher recording density, the advent of which is in 1 to 2 years. The heat-assisted recording system is a technique that enables magnetic recording in higher recording density by radiating a laser beam or the like onto the surface of the magnetic recording medium, heating the magnetic recording medium with light absorbed therefrom, reducing the coercivity of the magnetic recording layer using the resultant increased temperature, and assisting the head in writing signals. In this heat-assisted recording system, the protective layer and lubricating layer are expected to be heated up to approximately 300° C. or higher due to the heat resulting from the increased temperature of the magnetic recording layer. Therefore, the protective layer and the lubricating layer that are employed in such a heat-assisted recording system are required to be corrosion-resistant and durable to the heat from the laser radiation (referred to as “heating with laser beam”), as described above.\nUnfortunately, employing such a heat-assisted recording system in a conventional magnetic recording medium with a DLC protective layer to write signals, leads to quality degradation in the protective layer and lubricating layer due to the heat from the laser beam. This consequently causes an increase in the number of errors due to deterioration of the corrosion resistance of the protective layer and the lubricating layer, a shortening of the life of the magnetic recording medium due to degradation of the durability of the protective layer and the lubricating layer, and an increase in the risk of damage to the magnetic recording medium. As a measure to prevent these phenomena, Patent Document 1 discloses a method for forming a protective layer on a heat-insulating layer provided on a magnetic recording layer. Patent Document 1 describes that such materials of low thermal conductivity as SiO2, TiO2, and ZrO2 are preferred as the heat-insulating layer and that the heat-insulating layer is made by reactive sputtering method using O2 gas for forming a film with a target of each material.\nHowever, after executing this method, the inventors of the present application discovered that the O2 gas for forming the heat-insulating layer causes oxidation and hence deterioration of the magnetic recording layer and that the method is therefore not for practical use.\nPatent Document 1: Japanese Patent Application Publication No. 2010-153012"}
{"text": "The present disclosure relates to a support unit that supports a member to be supported, the support unit being removably mounted in an apparatus main body, and to an image forming apparatus incorporated with the support unit.\nIn an image forming apparatus, a toner image is formed on a photoconductor drum by a developing device, and the toner image is transferred onto a sheet in a transfer unit. The image forming apparatus further includes a fixing unit. The sheet on which the toner image has been transferred undergoes a fixing process by the fixing unit and is then discharged. Some of image forming apparatus thus configured to form an image on a sheet includes a support unit that supports a predetermined member, the support unit being removably mounted in the apparatus main body.\nThe support units thus far developed include a drum unit that supports the photoconductor drum, a developer unit that supports the developing device, and a process unit that integrally supports the photoconductor drum and the developing device. Such support units include a connector unit for electrical connection. When the support unit is mounted in the apparatus main body the connector of the apparatus main body and the connector unit of the support unit are electrically connected, so that driving power and control signals are supplied to the support unit from the apparatus main body. Some of such support units include a cover member that conceals the connector unit.\nHere, it is preferable that the support unit includes a grasp portion that facilitates an operator to mount or remove the support unit in and from the apparatus main body. For example, the grasp portion is provided on a cover member that opens and closes the apparatus main body."}
{"text": "Particles able to absorb or scatter light of well-defined colors have been used in applications involving detection, absorption, or scattering of light, for example medical diagnostic imaging. Such particles are typically colloidal metal particles. The term colloidal conventionally refers to the size of the particles, generally denoting particles having a size between about 1 nanometer and about 1 micron.\nSmall particles made from certain metals that are in the size range of colloidal metal particles tend to have a particularly strong interaction with light, termed a resonance, with a maximum at a well-defined wavelength. Such metals include gold, silver, platinum, and, to a lesser extent, others of the transition metals. Light at the resonance wavelength excites particular collective modes of electrons, termed plasma modes, in the metal. Hence the resonance is termed the plasmon resonance.\nBy selecting the metal material of a colloidal particle, it possible to vary the wavelength of the plasmon resonance. When the plasmon resonance involves the absorption of light, this gives a solution of absorbing particles a well-defined color, since color depends on the wavelength of light that is absorbed. Solid gold colloidal particles have a characteristic absorption with a maximum at 500-530 nanometers, giving a solution of these particles a characteristic red color. The small variation in the wavelength results from a particle size dependence of the plasmon resonance. Alternatively, solid silver colloidal particles have a characteristic absorption with a maximum at 390-420 nanometers, giving a solution of these particles a characteristic yellow color.\nUsing small particles of various metals, particles can be made that exhibit absorption or scattering of selected characteristic colors across a visible spectrum. However, a solid metal colloidal particle absorbing in the infrared is not known. Optical extinction, in particular absorption or scattering, in the infrared is desirable for imaging methods that operate in the infrared. Further, optical communications, such as long distance phone service that is transmitted over optical fibers, operate in the infrared.\nIt has been speculated since the 1950's that it would be theoretically possible to shift the plasmon resonance of a metal to longer wavelengths by forming a shell of that metal around a core particle made of a different material. In particular, the full calculation of scattering from a sphere of arbitrary material was solved by Mie, as described in G. Mie, Ann. Phys. 24, 377 (1908). This solution was extended to concentric spheres of different materials, using simplifying assumptions regarding the dielectric properties of the materials, by Aden and Kerker, as described in A. L. Aden and M. Kerker, J. of Applied Physics, 22, 10, 1242 (1951). The wavelength of the plasmon resonance would depend on the ratio of the thickness of the metal shell to the size, such as diameter of a sphere, of the core. In this manner, the plasmon resonance would be geometrically tunable, such as by varying the thickness of the shell layer. A disadvantage of this approach was its reliance on bulk dielectric properties of the materials. Thus, thin metal shells, with a thickness less than the mean free path of electrons in the shell, were not described.\nDespite the theoretical speculation, early efforts to confirm tunability of the plasmon resonance were unsuccessful due to the inability to make a particle having a metal shell on a dielectric core with sufficient precision so as to have well-defined geometrical properties. In these earlier methods, it was difficult to achieve one or both of monodispersity of the dielectric core and a well-defined controllable thickness of a metal shell, both desirable properties for tuning the plasmon resonance. Thus, attempts to produce particles having a plasmon resonance in keeping with theoretical predictions tended to result instead in solutions of those particles having broad, ill-defined absorption spectra. In many instances this was because the methods of making the particles failed to produce smooth uniform metal shells. Rather, the methods tended to produce isotropic, non-uniform shells, for example shells having a bumpy surface.\nHowever, one of the present inventors co-developed a novel method of making metallized nanoparticles (particles with a size between about 1 nanometer and about 5 microns) that was successful in producing metal-coated particles having narrow well-defined spectra. Further, one of the present inventors co-discovered that improved agreement with theoretical modeling of the metallized nanoparticles resulted from the incorporation in the theory of a non-bulk, size-dependent value of the electron mean free path. That is, improved agreement with theory was achieved by developing an improved theory applicable to thin metallization layers. Thus, in the improved theory a dependence of the width of the plasmon resonance on the thickness of the metallization layer was described.\nParticles having at least one substantially uniform metallization layer have been termed metal nanoshells. Nanoshell structures that exhibit structural tunability of optical resonance's from the visible into the infrared can currently be fabricated. For example, complete nanoparticle metallization shell layers containing gold have been demonstrated. Gold has the advantage of a strong plasmon resonance that can be tuned by varying the thickness of the coating. More generally, the resonance may be tuned by varying either the core thickness or the thickness of the coating, in turn affecting the ratio of the thickness of the coating to the thickness of the core. This ratio determines the wavelength of the plasmon resonance. A further advantage of gold-coated particles is that they have shown promise as materials with advantages in imaging and diagnostics. In particular, they have utility as band-pass optical filters, impeding the photo-oxidation of conjugated polymers, and in conjunction with sensing devices based on surface enhanced Raman substrates.\nPresent methods for making nanoshells involve purifying suspensions of various intermediates, as well as purifying a suspension of the nanoshell products. Methods of making nanoshells are disclosed, for example, in U.S. Pat. No. 6,344,272 and in S. Oldenburg, R. D. Averitt, S. Westcott, and N. J. Halas, “Nanoengineering of Optical Resonances”, Chemical Physics Letters 288, 243-247 (1998), each hereby incorporated herein by reference. In particular, a method for making nanoshells may include coating linkers onto substrate particles so as to form linker-coated nanoparticles, seeding metal colloids onto the linkers so as to form colloid-seeded nanoparticles, and reducing metal onto the metal colloids so as to form nanoshells. Further, at various stages, the method includes purification of linker-coated nanoparticle suspensions, purification of colloid-seeded nanoparticle suspensions, and purification of nanoshell suspensions. Corresponding undesirable byproducts correspondingly include, for example, excess linkers, excess metal colloids, and excess metal ions, respectively.\nPresent methods for purifying nanoshells and intermediates thereof rely on a technique for purifying product or intermediate thereof that involves centrifugation and redispersal. The centrifugation and redispersal may be repeated a suitable number of times to achieve the desired level of purity. A disadvantage of this conventional laboratory method is that it is difficult to scale up to commercial scale. Thus, a scalable method of purifying nanoshells and intermediates thereof is desirable.\nReplacement of combined centrifugation and redispersal with filtration is known to those skilled in the art of purification. Filtration has the advantage that it is known to be a scalable method of purification, in particular filtration may be scaled up to commercial scale. However, previous laboratory scale attempts to use filtration to separate excess colloid from colloid-seeded nanoparticles serving as intermediates in making nanoshells have been unsuccessful. This lack of success occurred despite the typical size disparity between excess metal colloid, typically between about 1 nm and 5 nm, and colloid-seeded nanoparticles, typically between 50 nm and 5 μm. In particular, attempts by the present inventors to crossflow filter a suspension, as it results from the seeding reaction, of colloid-seeded nanoparticles and excess metal colloid using a filter having a 50 nm nominal pore size have demonstrated ineffective passage of the excess colloids through the filter in the filtrate. Rather, sufficient excess colloids were found to remain in the retentate with the colloid-seeded nanoparticles that subsequent reduction of metal so as to form metal nanoshells was unsuccessful.\nThus, notwithstanding the above-described teachings, there remains a need for a scalable method of making nanoshells, particularly when the method involves intermediate colloid-seeded nanoparticles."}
{"text": "This invention relates to a weapon sighting device that is not aligned with the barrel, but mounted off to one side of the weapon barrel and more particularly to a weapon sighting device for use by soldiers wearing chemical protective mask and clothing.\nSoldiers armed with the M16 series rifles while wearing their chemical protective mask and clothing, must engage in a series of contortion-like manipulations in order to attempt to sight a target with the standard M16 rifle sights. Current familiarization exercises now require soldiers to be able to engage a target at 25 meters while wearing the chemical protective mask and equipment. The current practice includes cocking the head to one side, tilting the rifle, estimating the impact point of the bullet, and other difficult exercises in order to permit sighting of the rifle while wearing the bulky chemical-protective mask and clothing. Many soldiers completing these exercises are confident in the ability of their chemical-protective equipment to keep them alive, but they develop sincere concerns over their abilities to defend themselves in such an environment.\nThe present invention involves an offset sighting device that allows one to aim the M16 weapon in a virtually normal position and effectively engage and properly sight targets in excess of 100 meters away. This greatly increases the effective range of the M16 series rifles under chemical protective mask and gear utilization conditions."}
{"text": "In order to extract ferromagnetic components which are contained in ores, the ore is ground into a powder and the powder obtained is mixed with water. A magnetic field generated by one or more magnets is applied to this suspension, as a result of which the ferromagnetic particles are attracted so that they can be separated from the suspension.\nDE 27 11 16 A discloses a device for separating ferromagnetic particles from a suspension, in which a drum consisting of iron rods is used. The iron rods are alternately magnetized during the rotation of the drum, so that the ferromagnetic particles adhere to the iron rods while other components of the suspension fall down between the iron rods.\nDE 26 51 137 A1 discloses a device for separating magnetic particles from an ore material, in which the suspension is fed through a tube which is surrounded by a magnetic coil. The ferromagnetic particles accumulate at the edge of the tube, while other particles are separated through a central tube which is located inside the tube.\nA magnetic separator is described in U.S. Pat. No. 4,921,597 B. The magnetic separator comprises a drum, on which a multiplicity of magnets are arranged. The drum is rotated oppositely to the flow direction of the suspension, so that ferromagnetic particles adhere to the drum and are separated from the suspension.\nA method for the continuous magnetic separation of suspensions is known from WO 02/07889 A2. This uses a rotatable drum in which a permanent magnet is fastened, in order to separate ferromagnetic particles from the suspension.\nIn known devices, a tubular reactor, through which the suspension flows, is used to separate the ferromagnetic particles from the suspension. One or more magnets are arranged on the outer wall of the reactor and attract the ferromagnetic particles contained in it. Under the effect of the magnetic field generated by the magnets, the ferromagnetic particles migrate onto the reactor wall and are held by the magnet arranged on the outside of the reactor. Although this allows effective separation, the separation method can however only be carried out discontinuously since after a particular quantity of the ferromagnetic particles have accumulated, the reactor has to be opened and the ferromagnetic particles removed. Only then is it possible for a new suspension to be supplied, or for the suspension already used once to be subjected to the separation method again."}
{"text": "Wind power is considered one of the cleanest, most environmentally friendly energy sources presently available, and wind turbines have gained increased attention in this regard. A modern wind turbine typically includes a nacelle fixed atop a tower, a generator and a gearbox housed with the nacelle, and a rotor configured with the nacelle having a rotatable hub with one or more rotor blades. The rotor blades capture kinetic energy of wind using known airfoil principles. The rotor blades transmit the kinetic energy in the form of rotational energy so as to turn a shaft coupling the rotor blades to a gearbox, or if a gearbox is not used, directly to the generator. The generator then converts the mechanical energy to electrical energy that may be deployed to a utility grid.\nAt least some known nacelles include a yaw system for controlling a perspective of the rotor relative to a direction of wind. Such wind turbines also include sensors for sensing a direction of the wind. Thus, the wind turbine controller is configured to adjust the yaw of the wind turbine via the yaw system based on the sensed wind direction.\nIf the wind turbine controller is offline, however, then there is no way to operate the yaw system. Without the yaw system, the wind turbine may be subjected to increased loads (e.g., asymmetric loads) that result from yaw misalignment which may contribute to significant fatigue cycles on the wind turbine components. As the wind turbine components become worn, the wind turbine becomes less effective.\nThus, there is a need for a new and improved system and method for wind turbine yaw control that addresses the aforementioned issues. More specifically, an autonomous system and method for controlling the yaw of the wind turbine that does not rely of the main controller of the wind turbine would be advantageous."}
{"text": "1. Technical Field\nThe subject matter described herein relates generally to online services and, more particularly, to data access lockdown.\n2. Related Background\nA typical person may have at least one online account. Commonly, people have online accounts that are associated with electronic communication (e.g., email), photo sharing services, video sharing websites, social networks, etc.\nA user of the online account shares information with other users. The user may directly share information with another user (e.g., in an email message). The user may indirectly share information with another user through, for example, an online social network (e.g., posting content on an online social network that is not designated for any specific user, but for a plurality of users having a common trait or characteristic).\nA user may share information with another user in situations where sharing may not be intended. For example, a user may store information (e.g., text, audio, video, or image file) in a shared storage system (e.g., cloud storage) but may not restrict other users from listening to, downloading, or otherwise accessing the information.\nAt one point or another, sharing information online may pose problems that need to be resolved, such as unauthorized access or data leakage."}
{"text": "There are a variety of known devices for connecting various parts of an orthopedic implant together. Among these are included connectors featured with TSRH-3D products of Medtronic Sofamor Danek, Inc. Such devices commonly include a part for holding an elongated member, such as a spinal rod, and a second member for holding a bone fixation element, such as a pedicle screw or hook. The two parts are connected together, preferably with the ability to rotate with respect to each other. Such devices are shown, for example, in U.S. Pat. No. 6,183,473 to Ashman and U.S. Pat. No. 5,643,263 to Simonson."}
{"text": "The increase of data and video mobile wireless traffic has increased the demand on wireless communication system's backhaul transport. One way to increase the capability of a wireless communication system's backhaul transport is to utilize an optical fiber transport network. However, creating an optical fiber network includes high construction costs including cable rights-of-way and fiber cable construction, not to mention third part fiber based transport leases. Thus, the owners of wireless communication systems seek to implement low cost bandwidth solutions.\nOverview\nGenerally described, a method of operating a control database for a wireless communication system, wherein wireless users and a wireless access hub will exchange user data over first wireless communication signals, wherein the wireless access hub and a wireless collector hub will exchange the user data over second wireless communication signals, and wherein the wireless collector hub and a core communication network will exchange the user data over third wireless communication signals. In an example, the method comprises receiving first geographic coordinates for the wireless access hub and an identity of the wireless collector hub and, in response, retrieving second geographic coordinates for the wireless collector hub, identifying a frequency band for the third wireless communication signals between the core communication network and the wireless collector hub, and identifying reserved time slots in the frequency band at the wireless collector hub.\nAdditionally, the method also comprises selecting available time slots in the frequency band at the wireless collector hub for the second wireless communications signals between the wireless collector hub and the wireless access hub, processing the first geographic coordinates and the second geographic coordinates to determine a first azimuth and a first angle for the wireless access hub and a second azimuth and a second angle for the wireless collector hub to optimize signal strength for the second wireless communication signals between the wireless access hub and the wireless collector hub, determining a power for the second wireless communications signals between the wireless collector hub and the wireless access hub based on a distance between the wireless collector hub and the wireless access hub, receiving a first query for the wireless access hub, and in response, transferring a first response indicating the first azimuth, the first angle, the frequency band, the available time slots, and the power, and receiving a second query for the wireless collector hub, and in response, transferring a second response indicating the second azimuth, the second angle, the frequency band, the available time slots, and the power. A wireless communication system for operating a control database is also described."}
{"text": "1. Field of the Invention\nThis invention relates to the field of chemical engineering and to solute-solvent separation and more particularly, to apparatus and method suitable for consumer and industrial use in distilling of water.\n2. Description of the Prior Art\nTechniques for the separation of a solute from a relatively non-volatile solvent are well known and have been used for a number of years. However, such techniques are often complex in operation, inexpensive and consume large amounts of energy. Thus, the application of these techniques, on a large scale, is limited. For example, the purification of water from a saline source or brackish supply or from tap water, requires a considerable amount of structure. Various approaches for water purification from sources of the type described include multi-stage flash evaporation making use of reduced pressure for evaporation; reverse osmosis using hydrostatic pressures; vacuum freezing; distillation (mostly batch type) based on the boiling of water; moisture recovery based on large amounts of air-water contact; and electrodialysis based upon ion exchanges.\nThe problems which the foregoing approaches encounter are significant, such as the fabrication of special membranes, and the of special alloys for protecting machinery against corrosion or scaling. Many of these techniques also require high energy consumption for mechanical or evaporation purposes. Similar considerations are often involved in the purification concentration and separation techniques for solution of relatively non-volatile solvents or solutes.\nIt is especially desirable to have available at a point of use, purified water both for drinking and other purposes. Firstly, the health benefits are well known; secondly, a centralized purified water supply could not maintain its high quality if it is delivered through conventional plumbing structures; thirdly, while there are water purification devices which are commercially available, they are not modestly priced, and they require frequent maintenance or replacement of key components while consuming large amounts of energy during operation.\nBecause of the numerous problems associated with conventional structures mentioned above, it is advantageous for household and public use to be able to conveniently and economically purify water from water supplies at the point of use of the water. The present invention relates in particular to improvements in distillation apparatus and methods which will satisfy the need for such a water source."}
{"text": "1. Field of the Invention\nThe present invention relates to an addressing unit, and more particularly to an addressing unit suitable for a storage device which reads and writes information by using an X address (row address) and a Y address (column address).\n2. Description of the Related Art\nVarious storage devices for storing information have been proposed. Among them, a storage device which reads and writes information by using the X address and the Y address has been known. FIG. 4 shows a prior art example of an addressing unit which used the X address and the Y address.\nFIG. 4 shows an addressing unit used in a storage device represented by a DRAM. Namely, it shows a so-called address multiplex type addressing unit in which an address signal comprising multiplexed X address and Y address is inputted from an address input terminal common to the row and the column in a plurality of times in a time division fashion in order to reduce the number of address input terminals to improve a packaging density on a printed circuit board.\nIn FIG. 4, numeral 41 denotes a buffer which translates an address signal ad0 inputted from an address input terminal 8 to an internal X address signal ad1 to be supplied to an X decoder 6 and an internal Y address signal ad2 to be supplied to a Y decoder 7. By translation process, the address signal ad0 inputted from the address input terminal 8 is delayed by a certain period of time.\nNumeral 2 denotes an internal control signal generation circuit which generates a control signal (X latch signal) .phi.X0 for taking the X address and a control signal (Y latch signal) .phi.Y0 for taking the Y address.\nWhen the X address is to be taken, the internal control signal generation circuit 2 supplies the X latch signal .phi.X0 to an X address latch 43. Thus, the internal X address signal ad1 outputted from the buffer 41 is latched in the X address latch 43. When the Y address is to be taken, the Y latch signal .phi.Y0 is supplied to a Y address latch 44. Thus, the internal Y address signal ad2 outputted from the buffer 41 is latched in the Y address latch 44.\nNumeral 45 denotes an X pre-decoder which temporarily decodes the internal X address signal ad1 latched in the X address latch 43, two bits at a time, before it is input to the X decoder. Namely, it decodes the input n-bit address signal to a signal which identifies 2.sup.n address lines. For example, when the internal X address signal ad1 latched in the X address latch 43 is a 3-bit signal, the X pre-decoder 45 outputs an X pre-decoded signal ad3 for identifying eight addresses.\nNumeral 46 denotes a Y pre-decoder which temporarily decodes the internal address signal ad2 latched in the Y address latch 44, two bits at a time, before it is input to the Y decoder 7. For example, when the internal T address signal ad2 latched in the Y address latch 44 is a 3-bit signal, a Y pre-decoded signal ad4 which identifies eight addresses is outputted from the Y pre-decoder 46.\nThe X decoder 6 generates and outputs an X select signal adX for selecting a word line of a memory cell (not shown) based on the X pre-decoded signal ad3 supplied from the X pre-decoder 45. The Y decoder 7 generates and outputs a Y select signal adY for selecting a bit line of the memory cell (not shown) based on the Y pre-decoded signal ad4 supplied from the Y pre-decoder 46.\nIn the case where the pre-decoders 45 and 46 are provided in front of the decoders 6 and 7, the number of transistors of the decoders 6 and 7 can be reduced and the process may be conducted faster in comparison with the case where the internal address signals ad1 and ad2 generated by the buffer 41 are directly inputted to the decoders 6 and 7.\nFIG. 5 shows another prior art example of the addressing unit which uses the X address and the Y address. FIG. 5 shows a synchronous addressing unit. The blocks, which are identical to those of the addressing unit as shown in FIG. 4 are designated by the identical numerals. Namely, in the addressing unit in FIG. 5, a burst counter 51 is used instead of the Y address latch 44 in FIG. 4. The burst counter 51 comprises a binary counter.\nIn the example of FIG. 4, a plurality of address data ad2 included in the internal Y address signal is sequentially supplied from the buffer 41 to the Y address latch 44 and latched therein, while in the example of FIG. 5, the first address data ad2' included in the internal Y address signal is latched in the burst counter 51 and then the burst counter 51 is counted up based on the address data so that all address data included in the internal Y address signal are sequentially generated.\nAs shown in FIG. 5, when the addressing unit is constructed synchronously by using the burst counter 51, the speed of the process is further increased than the addressing unit as shown in FIG. 4.\nFIG. 6 is a time chart showing an operation to latch the internal X address signal ad1 in the X address latch 43. In FIG. 6, the address signal ad0 (Valid-Data) is kept to be a high level or a low level in a set-up period tS0 and a hold period tH0 in synchronism with a falling edge of the control signal .phi.0 (FIGS. 6 (a) and (b)).\nThe address signal ad0 is translated to the internal X address signal ad1 by the buffer 41 and latched in the X address latch 43. The latching is conducted synchronous with the falling edge of the X latch signal .phi.X0 supplied from the internal control signal generation circuit 2 to the X address latch 43 (FIGS. 6 (c) and (d)). The X latch signal .phi.X0 is a signal generated through the processing of the control signal .phi.0 by the internal control signal generation circuit 2, and a certain delay time dt0 is caused for a time period from the input of the control signal .phi.0 to the acquirement of the X latch signal .phi.X0.\nUnder such a condition, in order to latch the internal X address signal ad1 in the X address latch 43 in synchronism with the falling edge of the X latch signal .phi.X0, namely, in order to keep the value of the internal X address signal ad1 (Valid-data) in the set-up time tS1 and the hold time tH1 to be a high level or a low level in synchronism with the falling edge of the X latch signal .phi.X0, it is necessary to match the delay time dt1 when the internal X address signal ad1 is generated in the buffer 41 from the address signal ad0, to the delay time dt0 in the internal control signal generation circuit 2.\nHowever, the delay time caused in the generation of the internal X address signal ad1 in the buffer 41 is generally shorter than the delay time dt0 caused in the generation of the X latch signal .phi.X0 in the internal control signal generation circuit 2. Accordingly, in the prior art, it is necessary to set the delay time in the buffer 41 wastefully long in order to latch the internal X address signal ad1 in the X address latch 43 in synchronism with the falling edge of the X latch signal .phi.X0.\nNamely, in spite of the fact that the process to translate the address signal ad0 to the internal X address signal ad1 in the buffer 41 is completed and it is ready to latch it in the X address latch 43, the X latch signal .phi.X0 is not yet generated at that time and the delay of the buffer 41 is wastefully long. This renders the implementation of high speed processing difficult. This problem is serious in the addressing unit of FIG. 4 as well as the addressing unit of FIG. 5 which aims the further speed-up of the process."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the determination of additives in metal plating solutions, and more specifically, to a method and system for analysis of additives in electrochemical plating solutions using a flow management system that minimizes loss of plating solutions and decreases sampling time.\n2. Description of the Related Art\nThe microelectronic manufacturing industry applies a wide range of thin film materials to form microelectronic structures. These thin film materials include metals and metal alloys such as, for example, nickel, tungsten, tantalum, solder platinum, copper, copper-zinc, as well as dielectric materials, including metal oxides, semiconductor oxides, and perovskite materials.\nA wide range of processing techniques has been used to deposit such thin films, including chemical vapor deposition, electroplating and electroless plating. Of these techniques, electrochemical processing techniques, i.e., electroplating, and electroless plating are the most economical.\nElectroless plating uses an electrolyte wherein suitable salts of the desired deposition metals are dissolved therein. The electrolyte can be water-based, organic-based or a molten salt, depending on the specific metal to be deposited.\nMetal deposition from solution has been used for depositing copper during the production process of printed circuit boards in the microelectronics industry. Electrochemical deposition is the preferred method because of several advantages over chemical or physical vapor deposition that include the avoidance of using a vacuum system and overall reduced costs. However, when electrochemical techniques are employed, optimum process performance cannot be obtained without careful control of the plating bath chemistry. Notably, the quality and efficiency of the electrochemical plating process depend on the bath composition, pH and temperature.\nCurrently, most semiconductor manufacturers replenish plating baths based on time or energy consumption (e.g., amp hours). However, the use of such replenishment methods is based on the fact that the plating system always functions under ideal conditions, when in fact, human error and system malfunctions can affect the depletion rate of the components. To overcome the uncertainty of timed replenishing operations, in-line monitoring systems have been adapted for use in plating systems that monitor the metal levels in the plating baths and replenish the baths when needed. These systems minimize contamination, however, the sampling process can be time-consuming and any delay in determining analyte concentrations in the solutions can produce problems in maintaining the desired properties of the electrodeposited metals. Furthermore, these in-line sampling techniques can consume a considerable amount of sample thereby wasting expensive plating solutions and/or adversely affecting the economics of the semiconductor manufacturing operation.\nAccordingly, there is a need in the art for a sampling method and system that overcome the shortcomings of the prior art, such as excessive wastage of plating solutions, slow reaction times from sampling to replenishing of plating baths, and contamination from previous testing samples."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may not constitute prior art.\nA common field failure with heavier portable power tools, such as portable saws, is a separation of the power cord from the tool due to an impulse load, or jerk, applied to the cord. This can occur when the tool is dropped while the plug end of the power cord is secured, or when a user carries the tool or lowers it from floor to floor or down a ladder by holding the power cord.\nTo isolate the power cord conductors or connections from the high forces imposed by jerking the power cord, the power cord according to the present disclosure is installed in the tool housing with a small service loop, or extra length of cable, between the cord clamp and the portion of the tool housing that secures the cord protector. A crimp-on device is installed on the power cord cable next to the cord protector. When the cord is subjected to jerking, the cable moves axially relative to the cord protector. As the cable moves, the crimp-on device compresses the extended end of the cord protector absorbing energy and reducing the forces transmitted to the cord set conductors or connections that are disposed within the housing.\nFurther areas of applicability will become apparent from the description provided herein. It should be understood that the description and specific examples are intended for purposes of illustration only and are not intended to limit the scope of the present disclosure."}
{"text": "The present invention relates to an optical system used to record data to and/or reproduce data from one or more types of optical discs having different thicknesses of cover layers and/or different data recording densities.\nThere are various types of optical discs such as a CD (compact disc) and a DVD (digital versatile disc). The DVD has the cover layer thinner than that of the CD, and has a data recording density higher than that of the CD. For supporting both of the CD and DVD, the optical system for optical discs (i.e. an optical pick-up) is required to correct a spherical aberration which changes depending on the thickness of the cover layer of an optical disc being used and to change NA (numerical aperture) of a light beam so as to attain an effective beam diameter suitable for recording/reproducing operation of the optical disc being used.\nIn general, the optical system for optical discs has a light source for emitting the light beam, a coupling lens, and an objective lens. The coupling lens has the function of collimating the light beam emitted by the light source, or changing the degree of divergence of the light beam emitted by the light source. That is, the coupling lens is employed in the optical system to suppress aberrations or to enhance the efficiency of use of light.\nThere is a demand for decreasing the number of optical components in the optical system to decrease the cost of an optical pick-up and to further downsize the optical pick-up. Each of Japanese Patent Provisional Publications No. HEI 8-62496, HEI 8-334686, SHO 64-25113 and HEI 2-223906 discloses an optical pick-up which does not require the coupling lens. The optical system disclosed in each of the publications is configured to form a relatively large beam spot suitable for an optical disc having a relatively low recoding density and having a relatively thick cover layer (e.g. CD or CD-R). Since the optical system has a relatively small NA, the relatively large beam spot is attained.\nIt is noted that aberrations are caused in the optical system disclosed in each of the publications if a position of the objective lens is changed with respect to a position of the light source. More specifically, a spherical aberration is caused when the objective lens is moved in a direction of its central axis (optical axis) for a focusing operation. A coma and astigmatism are caused when the objective lens is moved in a direction perpendicular to the optical axis from a reference axis for a tracking operation because in this case off-axis light is incident on the objective lens.\nAs used herein, the term reference axis represents an axis including the optical axis of the objective lens and passing through the light source in a condition in which the objective lens is not moved by the tracking operation in the direction perpendicular to the optical axis. Also, the term “reference position” of the objective lens is defined as an initial position of the tracking operation.\nIt is required that aberrations are sufficiently suppressed for performing the recording/reproducing operation with high accuracy. In particular, aberrations having an asymmetric property such as a coma and astigmatism are main factors that deteriorate quality of a recording (optical) signal and a reproducing (optical) signal. Therefore, it is particularly required that the optical system is configured to sufficiently suppress the aberrations having an asymmetric property.\nHereafter, the term “tracking shift movement” represents movement of the objective lens in the direction (i.e. a radial direction of an optical disc) perpendicular to the optical axis due to the tracking operation, and a term “focusing shift movements” represents movement of the objective lens in the direction of the optical axis due to the focusing operation.\nThe coma caused in a cover layer of the optical disc when the optical disc tilts relative to the optical axis of the objective lens has such a property that it changes depending on the thickness of the cover layer. For this reason, both surfaces of the objective lens provided in the optical system disclosed in each of the publications are configured to be aspherical surfaces which cancel the coma caused in the cover layer of the optical disc. That is, the objective lens is configured as an aplanatic lens.\nBy contrast, sufficient correction of another aberration having an asymmetric property, i.e. astigmatism, is very difficult even if the objective lens of which both surfaces are asymmetrical surfaces is employed in the optical system. Therefore, the optical system disclosed in each of the publications can be used only for an optical disc, such as a CD, having relatively large tolerance to aberrations. That is, the optical system is specialized for an optical disc drive used for the optical disc having a relatively large tolerance to aberrations.\nIn other words, the optical system can not be used for an optical disc (such as a DVD) which has a relatively large recording density and requires a relatively small beam spot diameter.\nRecently, optical systems which support both of the CD and DVD (hereafter, referred to as a CD/DVD-compatible optical system) have been provided. However, such a conventional CD/DVD-compatible optical system can not correct a coma which changes depending on the thickness of the cover layer of the optical disc. Further, astigmatism is not corrected in the CD/DVD compatible optical system as in the case of the optical system of each of the publications.\nIt is understood that the optical system disclosed in each of the publications can not be used as the CD/DVD-compatible optical system because the amount of a coma changes depending on the thickness of the cover layer of the optical disc."}
{"text": "The present invention relates to a novel process for producing vitamin A ester and especially to a process for producing all trans-vitamin A ester.\nVitamin A is a known compound. Its chemical formula is represented as following: \nVitamin A ester of formula (I) such as, acetate or palmitate are used in great quantities as medicines, food additives and feed additives, etc. A variety of methods for the synthesis of vitamin A and its derivatives have been proposed.\nAmong those methods, there are three routines relating to Wittig or Wittig-Homer reaction: \nThe disadvantage of routine 1 lies in the fact that the triphenylphosphine reactant required for the synthesis is relatively expensive and that the byproduct of the reaction, i.e., Ph3PO, is water insoluble, thus it is difficult to isolate the desired product. In routine 2, Wittig-Horner reaction is employed to overcome the shortages existing in routine 1, the byproduct phosphonate is soluble in water, and thus can easily be isolated from desired product. However, the expected product is retinoic acid ester rather than vitamin A ester. In routine 3, intermediate aldehyde (II) in routine 1 and intermediate phosphonate compound (III) in routine 2 were employed to synthesize vitamin A ester directly.\nThe phosphonate compound (III) can be obtained by the following reactions: (see U.S. Pat. No. 4,916,4250) \nStarting from routine 3, it is an object of this invention to provide an improved process for more conveniently producing vitamin A ester, especially all trans-vitamin A ester in good yields. The advantages of this invention will become apparent from the following descriptions.\nThis object can be reached by a process for preparation of vitamin A ester, especially all trans-vitamin A ester of formula (I) which comprises treating compound represented by the formula (IV) with aldehyde represented by the formula (II) in an organic solvent in the presence of a base. \nwherein R1=3-methyl-1,3-pentadienyl or 3-methyl-1, 4-pentadienyl, R2 and R3 are identical or different and are C1-C4-alkyl\nR4 is C1-C15-hydrocarbyl.\nCompound (IV) includes:\nA; 3-methyl-5-(2, 6, 6-trimethyl-1-cyclohexen-1-yl)-1, 3-pentadienyl-phosphonic acid dialkyl ester \nB: 3-methyl-5(2, 6, 6-trimethyl-1-cyclohexen-1-yl)-1, 4-pentadienyl-phosphonic acid dialkyl ester or \nC: the mixture of A and B in any ratios\nIn CN 1097414A, compound (III), which was employed as starting material to synthesize vitamin A ester, was obtained by isomerization of compound (IV). However, according to the present invention, compound (IV) is directly employed as starting material to synthesize vitamin A ester. Due to the omission of isomerization step, vitamin A ester can be prepared in a more simply manner and higher yield according to the process of tie present invention,\nIn present invention, R2 and R3 are identical or different and represent C1-C4 alkyl, for example methyl, ethyl, n-propyl, isopropyl, n-butyl, iso-butyl, sec-butyl and tert-butyl, preferably methyl, ethyl, isopropyl. R4 represents saturated or unsaturated C1-C15 hydrocarbyl, for example, alkyl, alkenyl etc., preferably C1-C15 alkyl, for example, C1-C4-alkyl as mentioned above, pentyl, hexyl, heptyl, octyl, nonyl, decyl, undecyl, dodecyl, tridecyl, tetradecyl, pentadecyl and the isomer thereof, especially preferably methyl, ethyl, pentadecyl.\nIt is believed that the mechanism of this reaction is as follows: \nAccording to the process of this invention, vitamin A ester can be obtained by the reaction of compound (II) with intermediate carbanion (V), which is generated when compound (IV) is treated with a proper base. Suitable base is inorganic base or organic base. Preferred inorganic base includes, for example, alkali metal hydrides such as sodium hydride and potassium hydride. Preferred organic base includes organic base of alkali metal, for example, alkali metal alkoxides such as sodium tert-butoxide alkali potassium tert-butoxide, alkali metal salt of sulfoxides such as dimsyl sodium and dimsyl potassium, organolithium compounds such as methyllithium and n-butyllithium; organomagnesium halides such as ethylmagnesium chloride.\nThe amount of the base is not critical, the molar ratio of base to the compound of formula (IV) is generally from about 1 to about 3, preferably from about 1 to about 2. The molar ratio of the compound of formula (II) to the compound of formula (IV) is generally from about 1 to about 2.\nIn general, the reaction is carried out in an organic solvent. Any organic solvents which do not adversely affect the reaction can be employed. Preferred organic solvents include two types of organic solvents: (1) nonpolar aprotic solvents, for example, hydrocarbons such as benzene, hexane, cyclohexane, toluene; ethers such as tetrahydrofuran, diisopropyl ether; (2) polar aprotic solvents, for example, sulfoxides such as DMSO, ketones such as acetone, nitrites such as acetonitrile, amides such as DMF, HMPT. These solvents may be used either singly or in combination. Preferably, two types of solvents are used in combination. The amount of the solvent is not critical. Generally, it is advantageous that from about 0.05 to about 1 mole, preferably from about 0.1 to about 0.5 mole of compound (IV) are employed per liter of solvent. The temperature at which the reaction is carried out may be varied in wide rang upon the type of the used base. Generally, the suitable temperature is from about xe2x88x9270xc2x0 C. to about 70xc2x0 C., preferably from about xe2x88x9270xc2x0 C. to about 0xc2x0 C. The reaction is preferably carried out in an atmosphere of an inert gas such as helium, nitrogen or argon The time of reaction is not critical, generally from 30 minutes to 5 hours, preferably from 2 to 4 hours. Preferably, compound (IV) and base are added in the solvent prior to compound (II).\nBy the aforesaid reaction, vitamin A ester predominantly containing all trans-vitamin A is formed in good yields from the compound of formula (IV).\nThe resulting vitamin A ester of formula (I) can be isolated and separated from the reaction mixture in a known manner. For example, water or an aqueous solution of ammonium chloride etc., is added to the reaction mixture, and the organic phase is separated from the mixture. If necessary, the organic phase is washed with water and dried over anhydrous sodium sulfate, and the solvent is evaporated under reduced pressure to separate vitamin A ester. If necessary, the product may be purified, for example, by recrystallization to give vitamin A ester of high purity.\nAccording to the present invention, all trans-vitamin A derivatives can be obtained in good yields and high purity, the content of other isomers such as 9-cis, 11-cis, and 13-cis isomers is very low. It is easy to obtain all trans-vitamin A derivatives from the formed product by conventional purifying means such as recrystallization,\nCompound of formula (IV) can easily be obtained according to known methods or in analogy to the known methods, for example, it can be obtained from starting material xcex2-ionone according to the method described in U.S. Pat. No. 4,916,250, \nAs to the synthesis of compound of formula (II), reference is made to, for example,\nReif W. et al., Chemie. Ing. Techn, 1973; 45 (10a); 648;\nEletti-Bianchi G. et al., J. Org. Chem. 1976; Vol. 41: 1648;\nJ. Org. Chem, Vol. 42, 1977, 2939;\nU.S. Pat. No. 5,527,952;\nJ. Org. Chem. Vol, 44, 1979, 1716"}
{"text": "1. Field of the Invention\nThe present invention relates to the configuration of two-dimensional code which can be detected from an image, and to a method and device for detecting two-dimensional code from an image.\n2. Description of the Related Art\nA two-dimensional code, in which a two-dimensional region is sectioned according to a certain rule and information is expressed by the brightness/darkness of each of the section regions, is known. Normally, such code is captured with a camera, the two-dimensional code is distinguished and detected from the background, and information is read out based on the brightness of each sectioned region appropriated therein. Various forms of two-dimensional code are known as conventional art, and the two-dimensional code shown in FIGS. 9 and 10 is a type thereof. This two-dimensional code will hereafter be referred to as “two-dimensional code 1”.\nThe two-dimensional code has an external shape of a square which is darker than a bright background, and the interior thereof starting from a predetermined spacing from the frame is divided into small squares, thus expressing bits by the brightness/darkness of the squares.\nAlso, in the field of mixed reality wherein a virtual space is displayed so as to be merged with real space, square markers disposed in the real space are used as a way to estimate the position and orientation of cameras in the real space. The demand for positioning multiple square markers in real space has brought about a proposal for square markers each of which have individual ID information within the square so as to enable each square to be identified uniquely. Examples of the method disclosed in X. Zhang, S. Fronz, N. Navab: “Visual Marker Detection and Decoding in AR systems: A Comparative Study,” Proc. of International Symposium on Mixed and Augmented Reality (ISMAR' 02) 2002 (hereinafter “Zhang et al.”) include the marker 1 shown in FIG. 12 and the marker 2 shown in FIG. 13.\nThe marker 1 is of a configuration wherein a large square of which the interior is sectioned into small squares, and several small squares along one side of the large square, are arranged. The structure of the interior of the large square is the same as that of the two-dimensional code 1 in the above-described related art, the interior being equally divided into small squares with a certain spacing between the outer frame of the square and the interior thereof. Bits can be represented by the brightness/darkness of each in the squares on the interior. The marker 1 also has a region of small squares on the outside of the large square, which is used for determining the reference side of the square, that is to say the direction of rotation of the square. Also, in addition to the direction of rotation, other types of information can be represented by the relative position between the small squares positioned outside as to the large square, and the number thereof.\nThe marker 2 is the same as described in Zhang et al. and the two-dimensional code 1 and marker 1 described above regarding the fact that the shape is a square which is darker than a bright background. The point that differs is the shape of the regions for assigning bits within the square; with this marker, circular regions are used. However, the fact remains that regions are appropriated as to the interior region of a square, with a certain spacing between the outer frame of the square and the interior thereof, albeit circles.\nThe two-dimensional code 1, marker 1, and marker 2, described as related art, are each based on the same technique of having an external shape of a square which is darker than a bright background, and the interior thereof starting from a predetermined spacing from the frame is equally divided into sections, thus expressing bits by the brightness/darkness of the sectioned regions.\nNow, the predetermined spacing between the square which is the outer shape, and the interior region thereof, is set so as to be the same as or larger than the size of each of the sectioned regions into which the interior has been sectioned for the purpose of carrying information. The case of the two-dimensional code 1, for example, will be described with reference to FIG. 11. The symbol d1 in FIG. 11 represents the spacing provided between the outer shape square and the interior bit region. On the other hand, d2 represents the length of one side of each of the interior bit squares. As can be understood from this drawing, d1 is set so as to be longer than d2. Also, the marker 1 will be described in the same way with reference to FIG. 14. As with the case above, d1 in FIG. 14 represents the spacing provided between the outer shape square and the interior bit region, and d2 represents the length of one side of each of the interior bit squares. In the case of the marker 1, d1 and d2 are equal. Hereinafter, the two-dimensional code 1, marker 1, and marker 2 will be collectively referred to as “two-dimensional code”.\nAs can be understood so far, with the two-dimensional code of the related art, the relationship between the spacing d1 between the outer shape square and the interior bit region, and the length d2 of one side of each of the interior bit squares, is d1≧d2.\nWith the two-dimensional code of the related art, there is a discrepancy between the minimum size of the two-dimensional code at which the outer shape can be recognized, and the minimum size of the two-dimensional code at which the inner bit region can be read. That is to say, the sizes d1 and d2 have not been optimized.\nAccordingly, there has been the problem that while the square outer shape of the two-dimensional code can be recognized, the interior bits could not be read out with a high degree of accuracy. From the opposite perspective, with the minimum size at which the interior bits can be read out, the space between the outer frame and the interior region that is necessary for recognizing the outer shape is excessively great, meaning that the area of the two-dimensional code in the image is too great."}
{"text": "1. Field of the Invention\nThe present invention relates to a dumbell and hand-grip exerciser in conjunction with a ball throwing and catching handgoal device.\n2. Description of Prior Art\nHeretofore, various exercise and ball game catching devices, and the like, have been developed. None of these, however, has the features of the present invention. Example;\n______________________________________ UNITED STATES PATENTS 4,629,186 12/1986 ALDRIDGE 194,294 12/1962 GREIF (design) 749,147 1/1904 AYELOTT 208,787 10/1967 PRESSCOTT (design) 4,632,384 12/1986 BRIGHT 1,026,515 5/1912 KORTH 4,756,522 7/1988 SANDOVAL 4,155,552 5/1979 JACOBO et al 1,175,035 3/1916 WOOSTER 3,711,096 1/1973 CRAMP et al 3,061,311 10/1962 VON ARHHERN 3,289,246 12/1966 DEYE 1,562,432 11/1925 BERSCH FOREIGN PATENTS 2,523,339 3/1967 DENKO FRANCE 1,023,996 2/1958 BUCHHOLZ GERMANY 29,241 9/1913 GODSON BRITAIN ______________________________________"}
{"text": "1. Field of the Invention\nThe present invention relates to a visual acuity testing apparatus, more particularly, the present invention pertains to a visual acuity testing apparatus of a subjective type used for inspecting visual acuity.\n2. Related Art Statement\nIn conventional subjective type visual acuity testing apparatuses, it is required to lengthen the distance between an eye to be tested and an apparatus for inspecting visual acuity in order to eliminate an influence due to the ability of accommodation of the eye and therefore, these apparatuses have been constructed so that a visual acuity table is placed at 6 meters distant from the eye to carry out the inspection of visual acuity.\nFor that reason, a quite wide space for carrying out the visual acuity inspection should be secured when using such a conventional subjective visual acuity testing apparatus. However, it is particularly difficult for opticians and ophthalmologists to secure or maintain such wide space since, in general, the floor area thereof is very small and limited.\nMoreover, in the conventional subjective visual acuity testing apparatuses, a support arm is vertically disposed on and fixed to the top of a table to support a unit for inspecting visual acuity at an upper portion of the support arm.\nWhen such a subjective visual acuity testing apparatus is used to inspect eyes, it is also necessary to secure a space to be occupied by a consultant other than the space for eye examination.\nHowever, it is quite difficult for opticians and ophthalmologists, the floor areas of which are quite small, to secure both areas for carrying out eye examination and consultation."}
{"text": "A) Field of the Invention\nThis invention relates to an electronic musical apparatus, and more in detail, an electronic musical apparatus that can display lyrics and a chord name on other electronic musical apparatus.\nB) Description of the Related Art\nIn an electronic musical apparatus that has an automatic musical performing function such as an electronic musical instrument, when music data including lyrics data is reproduced, it is well-known that an external displaying apparatus displays lyrics via a video-out device (image data output circuit), for example refer to JP-A 2002-258838.\nIn the above-described prior art, lyrics corresponding to music data are output to an external apparatus as image data, and lyrics can be displayed on a separated displaying device and a displaying device that has a large screen.\nIn the prior art, however, image data (image signals) for displaying lyrics is generated based on lyrics data and image data is transmitted to an external apparatus via a video-out device, this kind of apparatus is expensive since the video-out device is generally expensive."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming device such as a copying machine, a digital multi-functional device, a printer, and the like.\n2. Background Information\nImage forming devices such as copying machines, digital multi-functional devices, printers, and the like are now widely used in offices. When these types of image forming devices perform tasks such as copying original documents and printing data received by facsimile, noises will be generated thereby such as the sound of the transfer drum rotating and the sound of paper being transported, printed and discharged.\nBecause of that, it will be difficult to hear another person's voice over the telephone or during a face-to-face conversation if, for example, these conversations occur in the vicinity of an image forming device performing these types of operations.\nIn response to this problem, image forming devices that, for example, are able to somewhat prevent this noise from being generated, or that have high soundproofing capabilities so that noise generated during copying and printing will not leak out of the device, have recently been developed (see, for example, Japanese Published Patent Application No. H06-194885).\nHowever, it is difficult to completely prevent the aforementioned noise from being generated, and also difficult to completely prevent this noise from leaking out of the image forming device. In addition, the cost of the image forming device will increase if a configuration that improves the soundproofing capabilities thereof or suppresses the generation of noise therein is adopted.\nIn addition, sounds generated by the image forming device not only include the noise generated during printing, but also include, for example, notification sounds that inform a user that the operation keys displayed on the operation panel of the image forming device have been pressed, a notification sound that indicates receipt of facsimile by the image forming device, a notification sound that indicates that printing is completed, a notification sound that warns a user of the existence of a problem such as a lack of paper, and the like. Because these sounds are not very loud, there is almost no difficulty for a person to speak to someone over the telephone or in person while in the vicinity of the image forming device. However, there are times in which these notification sounds will be drowned out by the background noise of the office because they are not loud enough. Thus, when background noise is being generated in the office, there will be times when, for example, the user will not be able to recognize that an operation key has been pushed, will not notice that data has been received by facsimile, or will not notice that the device has run out of paper."}
{"text": "This invention relates to two-piece gate valves and, in particular, a two-piece gate valve having the gate pieces operatively spaced by a spacer ring.\nDouble disc gates have been proposed for gate valves because of certain inherent functional and manufacturing advantages. Such constructions, in view of the larger freedom of movement of the independent discs, can conform to larger manufacturing variations in the discs and the valve seats while, at the same time, provide complete sealing despite distortion of the seats due to differential heating and cooling rates and system stresses."}
{"text": "The present invention is directed to organic solvents and electrolytes for electrochemical cells, particularly lithium ion cells, and to electrochemical cells exhibiting good low temperature performance.\nState-of-the-art lithium ion cells typically include a carbon (e.g., coke or graphite) anode intercalated with lithium ions to form LixC; an electrolyte consisting of a lithium salt dissolved in one or more organic solvents; and a cathode made of an electrochemically active material, typically an insertion compound, such as LiCoO2. During cell discharge, lithium ions pass from the carbon anode, through the electrolyte to the cathode, where the ions are taken up with the simultaneous release of electrical energy. During cell recharge, lithium ions are transferred back to the anode, where they reintercalate into the carbon matrix.\nLithium ion rechargeable batteries have the demonstrated characteristics of high energy density, high voltage, and excellent cycle life, making them more attractive than competing systems such as Nixe2x80x94Cd and Nixe2x80x94H2 batteries. However, few state-of-the-art lithium ion cells perform well at low temperatures making them unsuitable for many terrestrial and extra-terrestrial applications. Many scheduled NASA missions demand good low temperature battery performancexe2x80x94without sacrificing such properties as light weight, high specific energy, long cycle life, and moderate cost. The Mars Exploration Program, for example, requires rechargeable batteries capable of delivering 300 cycles with high specific energy, and the ability to operate over a broad range of temperatures, including the extremely low temperatures encountered on and beneath the surface of Mars. Mars Rovers and Landers require batteries that can operate at temperatures as low as xe2x88x9240xc2x0 C. Mars Penetrators, which will penetrate deep into the Martian surface, require operation at temperatures lower than xe2x88x9260xc2x0 C.\nTo be used on the Mars missions and in low earth orbit (LEO) and geostationary earth orbit (GEO) satellites, as well as in terrestrial applications, lithium ion rechargeable batteries should exhibit high specific energy (60-80 Wh/Kg) and long cycle life (e.g.,  less than 500 cycles).\nUnfortunately, state-of-the-art lithium ion cells typically exhibit poor capacities below 0xc2x0 C. This is primarily due to limitations of the electrolyte solutions, which become very viscous and freeze at low temperatures, resulting in poor ionic conductivity. In addition, the surface film (i.e., solid electrolyte interphase, SEI) that forms on the electrodes either builds up over the course of repeated charge/discharge cycling or becomes highly resistive at lower temperatures. Ideally, the SEI layer on the carbon anode needs to be protective toward electrolyte reduction and yet conductive to lithium ions to facilitate lithium ion intercalation, even at low temperature.\nA number of factors can influence the low temperature performance of lithium ion cells, including (a) the physical properties of the electrolyte (i.e., conductivity, melting point, viscosity, etc.), (b) the electrode type, (c) the nature of the SEI layers that can form on the electrode surfaces, (d) cell design, (e) electrode thickness, separator porosity and separator wetting properties. Of these, the electrolyte properties have the predominant impact upon low temperature performance, as sufficient electrolyte conductivity is a necessary condition for good performance at low temperatures. Ideally, a good low temperature performance electrolyte should possess a combination of several critical properties, including high dielectric constant, low viscosity, adequate Lewis acid-base coordination behavior, as well as appropriate liquid ranges and salt solubilities in the medium.\nConventional electrolytes employed in state-of-the-art lithium ion cells have typically consisted of binary mixtures of organic solvents, for example high proportions of ethylene carbonate, propylene carbonate or dimethyl carbonate, within which is dispersed a lithium salt, such as LiPF6. Examples include 1.0M LiPF6 in a 50:50 mixture of ethylene carbonate/dimethyl carbonate, or ethylene carbonate/diethyl carbonate. Such electrolytes do not perform well at low temperature because they become highly viscous and/or freeze.\nIt can be seen, therefore, that a clear need exists for improved organic solvents, electrolytes, and electrochemical cells capable of performing well at low and moderate temperatures, with high specific energy and high cycle lives.\nThe present invention provides novel organic solvent systems, electrolytes, and electrochemical cells characterized by improved low temperature performance, including high conductivity, good cycle life, good discharge characteristics, good stability and self-discharge characteristics, and excellent compatibility with the cell components, as well as excellent room temperature and elevated temperature performance. Lithium ion cells containing the new electrolytes are ideal for use in portable electronic products, space vehicles, and other applications.\nIn one embodiment of the invention, an organic solvent system comprises a ternary mixture of ethylene carbonate (EC), dimethyl carbonate (DMC), and diethyl carbonate (DEC), preferably an equal volume mixture of each. Fluorinated analogsxe2x80x94especially perfluorinated analogsxe2x80x94can be used in place of one or more of the solvents.\nIn a second embodiment of the invention, an organic solvent system comprises a mixture of organic carbonates and at least one aliphatic ester, preferably an alkyl or fluoroalkyl ester. Thus, ternary, quaternary, and higher solvent systems are provided. One such solvent system comprises a ternary mixture of EC, DMC, and MA (methyl acetate). Another solvent system comprises a quaternary or higher mixture of EC, DMC, DEC and at least one alkyl or fluoroalkyl ester.\nIn still another embodiment, an improved organic solvent system includes two or more alkyl carbonates, preferably EC, DMC, DEC, and/or PC (propylene carbonate), and an asymmetric alkyl carbonate.\nIn yet another embodiment, the solvent system includes a compound having the formula LiOX (where X is xe2x80x94R, xe2x80x94COOR, or xe2x80x94COR, where R is alkyl or fluoroalkyl) or another basic species that can effectively catalyze the described disproportionation reactions. Fluorinated (especially perfluorinated) analogs of one or more of the co-solvents or additives can also be used. Where the organic solvent system includes an asymmetric alkyl carbonate, such as ethyl methyl carbonate, it can be added directly to the solvent system or generated in situ by including a lithium alkoxide or similar basic species in the mixture. Thus, the present invention also provides a method for making organic solvent systems for electrochemical cells in which asymmetric alkyl carbonates can be produced.\nThe addition of lithium methoxide or a related basic species also has the observed benefit of improving the SEI formation characteristics of carbonate-based (and other) electrolytes, which contributes to improved cell performance, especially at low temperature, due to low electrode polarization behavior. Lithium methoxide has previously been detected as a by-product formed from electrolyte reactions in lithium-ion cells which have been cycled or exposed to lithiated graphite. In fact, the addition of alkoxides and related basic compounds to lithium ion cell electrolytes has now been found to have a beneficial effect, irrespective of whether the additive(s) can facilitate disporportionation/exchange reactions resulting in asymmetric carbonate formation. The present invention, therefore, includes a method of making an improved electrolyte by adding to an electrolyte solvent system a small amount of a compound of the formula LiOX (as described herein), or a similar basic species.\nIn addition to novel organic solvent systems, the invention also provides a variety of improved electrolytes, comprising a lithium salt having high ionic mobility, dispersed in an organic solvent system. A preferred salt is LiPF6. The electrolytes are particularly well-suited for use in lithium ion cells, especially where low temperature performance is required.\nIn another embodiment of the invention, an improved, reversible electrochemical cell comprises an anode, a cathode, and an electrolyte as described herein. Preferably, the anode and cathode are of the insertion type, i.e., a coke, graphite, or modified carbon anode intercalated with lithium ions, and the cathode is a lithiated metal oxide, such as LiCoO2, or a similar material. In operation of the cell, lithium ions pass from the carbon through the electrolyte to the cathode, where they are taken up. During recharge, lithium ions are transferred back to the anode, where they reintercalate into the carbon matrix.\nLithium ion cells prepared in accordance with this invention have high specific energy, energy density, operating voltage, and coulombic efficiency, and low self-discharge tendencies. Cell performance at low temperatures ( less than xe2x88x9220xc2x0 C.)is notably improved over comparable, state-of-the-art systems. When used in upcoming space missions, the invention promises the additional benefits of improved reliability, and broader operating range."}
{"text": "1. Field of the Invention\nThis invention relates to a process and apparatus for enhancing the growth of biomass within a reaction vessel and, in particular, enhancing biomass growth in reaction vessels which handle effluent and sewage.\n2. Discussion of the Prior Art\nIn recent years, a process for treating industrial effluent and sewage involving the use of biomass support bodies has shown itself as an attractive alternative to the conventional activated sludge treatment process. In the conventional process, sewage or industrial effluent in a reaction vessel is aerated and agitated to stimulate the metabolism of a biological population within the vessel with various impurities in the effluent. The secondary sludge produced by this interaction is removed by sedimentation\nIn the biomass support body process, the reaction vessel is a vertical structure and it contains a number of free moving biomass support bodies Effluent and air are introduced at the base of the vessel and, as the effluent is oxygenated and propelled upward by the air, the biological population in the biomass support bodies reacts with the oxygenated effluent to produce carbon dioxide and additional biomass. This process is described in detail in U.S. Pat. No. 4,419,243 which also describes a suitable apparatus for accomplishing the process.\nThe biomass support body process has demonstrated that the capacity of the reaction vessel used in the process for a given through flow of effluent can be approximately one fifth of that of a conventional aeration tank and so the process offers commensurate savings in the cost of an effluent treatment plant. However, even more savings could be achieved if the critical variables of the biomass growth process, such as the pH level, control of recirculated solids and control of biomass populations, could be more precisely controlled. Additionally, it would be desirable to completely control the movement of the biomass support bodies through the vessel independent of the flow of effluent as opposed to relying on the upward flow of the effluent to suspend or fluidize the bodies, as is required in the existing biomass support processes. Furthermore, a wider choice of materials and construction of the support bodies could be available if there were no requirement that biomass support bodies be suspended or fluidized."}
{"text": "The present invention generally relates to a delayed curing resin composition. More specifically, the present invention relates to a resin composition that allows the resin and catalyst to be mixed for an extended period creating a long pot life composition which cures only when subjected to elevated temperature.\nGenerally, epoxy coatings are well known in the art and due to their exceptional durability and structural properties epoxy based protective coatings have gained commercial acceptance as protective and decorative coatings for use on a wide variety of materials. For example, epoxy based protective coatings represent one of the most widely used methods of corrosion control. They are used to provide long term protection of steel, concrete, aluminum and other structures under a broad range of corrosive conditions, extending from atmospheric exposure to full immersion in highly corrosive environments. Further, epoxy coatings are readily available and are easily applied by a variety of methods including spraying, rolling and brushing. They adhere well to steel, concrete and other substrates, have low moisture vapor transmission rates and act as barriers to water, chloride and sulfate ion ingress, provide excellent corrosion protection under a variety of atmospheric exposure conditions and have good resistance to many chemicals and solvents. As a result, numerous industries including maintenance, marine, construction, architectural, aircraft and product finishing have adopted broad usage of epoxy coating materials.\nThe most common material utilized in the epoxy coating industry today is a multi-part epoxy material. In general the epoxy includes a first base resin matrix formed from a bisphenol material such as bisphenol A (BPA) and at least a second catalyst or hardener, although other components such as a pigment agent or an aggregate component may also be added. While the two parts remain separate, they remain in liquid form. After the two parts are mixed together, they begin a curing process that is typically triggered by exposure to heat, humidity or a ultra-violet light source, whereby the mixed material quickly begins to solidify. As a result, it is necessary to mix only a sufficient amount of compound such that it can be worked effectively before set up occurs. Accordingly, the use and application of these compounds is a tedious, slow and expensive proposition.\nThe hardeners are typically nitrogen-containing bases that are well known to the person skilled in the art as curing agents or curing accelerators for epoxy resins. Such systems have, however, only limited storage stability because those bases react with epoxides even at relatively low temperature, in some cases even at room temperature, which is manifested in an increase in the viscosity of the epoxy resin formulation and, on prolonged storage, results in hardening of the mixture. The greater the reactivity of the nitrogen-containing base, the lower the storage stability of the epoxy resin mixture and the shorter the pot life. For that reason, such systems are formulated as two-component systems, that is to say the epoxy resin and the nitrogen-containing base are stored separately and mixed only shortly before processing.\nThere has been no shortage of attempts at improving the storage stability of such systems by developing appropriate curing systems. The problem posed is particularly complex because, at the same time as the requirement for a high storage stability and a long pot life, there should not be any deterioration either in the reactivity at the desired curing temperature or in the properties of the fully cured materials. For example, adsorption techniques have been used to control and modify various types of chemical reactions. These techniques usually involve adsorbing a chemical reagent in an adsorbent material. Commonly used adsorbent materials for this purpose are materials having internal pore structure and active pore sites, and can consist of silica gel, certain types of carbon black, activated charcoal, and the like. In practice, when using an adsorbed chemical reagent in a process involving a controlled chemical reaction, the adsorbed chemical reagent and adsorbent is admixed with the reacting component at relatively low temperatures and subsequently heated to desorb the adsorbed component. Heating the adsorbent and adsorbate product desorbs the adsorbate reagent reactant making it available for a reaction with a reacting component. The mixture prior to being activated is relatively inert and fairly safe to handle.\nWhile the aforementioned reagent adsorption solves many of the problems in regard to process control, handling etc., there is in most instances a slow escape of the chemical reagent from the adsorbent. In many instances, this slow escape of chemical reagent from the adsorbent creates problems. The desorbed chemical reagent if in a reactive environment or reactive medium will allow a slow reaction between reagents to proceed. If the rate of escape of reactant is large and the resulting reaction is exothermic there is a possibility that the exothermic heat effect will generate sufficient heat in the mixture to desorb and activate the entire mixture, or at least accelerate desorption. Further, the slow escape of chemical reagent will cause product deterioration, and shortened shelf life of an adsorbed component mixture. In general, depending on the type of adsorbent and adsorbate chemical reagent, the rate of escape of adsorbate from adsorbent will vary. Even a small escape of adsorbate is objectionable and in certain instances can cause very serious effects.\nIn the reagent adsorption techniques known to the prior art, adsorbed chemical reagent when present in a surrounding medium containing a reactive medium, is not rendered completely inert. In general, an adsorbent has an open pore structure. A portion of the adsorbed chemical reagent is in immediate contact with the reactive medium and is therefore in a potentially reactive position. The adsorbed chemical reagent molecules, even though attracted and held in the active pore sites by Van der Waals forces will often be dislodged from the adsorbent by the normal molecular vibration of the chemical components, and will be free to react with the reactive medium. The tendency to dislodge the adsorbed adsorbate and the seriousness of this effect will vary with the type of adsorbate, the adsorbent and the reactive medium and the other possible components having a tendency to displace the adsorbate. Normally, the function of the adsorbent is to prevent or delay a reaction between the adsorbed chemical reagent and a reactive surrounding medium.\nThe more efficiently this function is performed, in general, the more desirable is the adsorption system. Therefore, adsorption of chemical reagents known to the prior art will not produce complete inertness, of a chemical reagent. Further, if the adsorbent in the chemical reagent combination is selected so that the adsorbate is very securely attached to the adsorbent, thus producing a very inert adsorbent adsorbate combination, it may require an extremely powerful displacing agent or heat effect to activate the material. For example, when the adsorbent adsorbate must be heated to extremely high temperatures in order to desorb the adsorbed reagent, other reagents in the mixture may be decomposed. This effect may completely prohibit the use of an adsorbed chemical reagent. An example is a decomposable blowing agent in a foamable mixture that will cure at a relatively low temperature. If the mixture is heated high enough to desorb the blowing agent, the reaction may proceed too rapidly and scorch, burn or cure the resin poorly.\nIn view of the foregoing, there is a need for a delayed curing epoxy resin composition. Further, there is a need for a delayed curing epoxy resin wherein the hardener and resin components can be fully blended yet the curing reaction still be delayed to provide the composition with a long pot life. There is still a further need for method of lining a pipe system whereby a liner is fully wet out with a blended two part epoxy composition yet the curing reaction is delayed for an extended period allowing the wet out liner to be stored and installed before the reaction is activated."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical disc apparatus such as a compact disc apparatus (CD), CD-ROM or other optical disc apparatus, or magneto-optic disc apparatus.\nFurther, the present invention relates to a tracking error signal calculation circuit used in the optical disc apparatus.\n2. Description of the Related Art\nAn optical pick-up is used for recording data or reading recorded data along a track (guide groove) of a disc recording medium of an optical disc apparatus. In the optical pick-up, optical parts such as a semiconductor laser, photodiode (PD), prism, and object lens are mounted.\nAt the time of recording or reading data to or from the disc recording medium, in order to eliminate the influence of the face vibration and wobbling of the track of the disc recording medium, and the face vibration and wobbling of the track by the tilting of the turntable of the disc drive, and precession motion, and the like, focusing servo control and track servo control are carried out.\nIn focusing servo control, the object lens for focusing (exactly focusing) the laser beam emitted from the semiconductor laser on the recording surface of the disc recording medium is positioned with respect to the surface of the disc recording medium.\nIn tracking servo control, the optical pick-up and object lens is positioned in the radial direction of the disc recording medium to position the laser beam emitted from the semiconductor laser at a desired track of the disc recording medium (to place it on-track).\nA focus error signal is used for the focusing servo control, while a tracking error signal is used for the tracking servo control.\nUsually, signals detected by two photodiodes are computed by the push-pull system to calculate the tracking error signal.\nAn offset will appear in the tracking error signal due to the push-pull system. When there is an offset, even if the tracking error signal indicates 0, if the tracking servo control is carried out by using that tracking error signal, the beam of light of the semiconductor laser will deviate from the center of the track, so it suffers from the disadvantage of poor control when performing the tracking control.\nAs factors behind the appearance of the offset in the tracking error signal, there are deviation of the optical axis of the object lens, tilting of the disc recording medium in the radial direction, an unbalance in groove shape on the disc recording medium, etc.\nAs the method for reducing the offset due to the above factors, a variety of countermesures have been made. For example, refer to \"Optical Disc Technology\", supervised by Morio Onoe, Radio Gijutsusha, page 91 to page 98.\nHowever, the inventors of the present application found that just the above-mentioned countermeasures for reducing the offset were not sufficient and that offset still existed in the tracking error signal. Accordingly, when using a tracking error signal calculated by the conventional push-pull system, the problem is encountered that accurate and stable tracking servo control cannot be carried out in the optical disc apparatus.\nFurther, if a deviation occurs in the balance of the tracking error signal accompanying movement of the field of vision of the object lens, the disadvantage is encountered that tracking servo control cannot be stably carried out."}
{"text": "(1) Field of the Invention\nThe present invention relates generally to exercise devices and, more particularly, to compact low impact therapeutic exercise devices designed for increasing circulation in a user.\n(2) Description of the Prior Art\nTypically, a person does not or cannot exercise while they are sitting. It does not matter whether one is sitting at a desk, or traveling in a plane, train, or automobile; one does not traditionally exercise while sitting. Additionally, many people spend a majority of their time in a sitting or similar position, typically engaged in one of the aforementioned situations. Being able to exercise while sitting in a locale that one does not traditionally exercise at would be beneficial to a person.\nThe benefits that one could receive from being able to exercise while sitting in a traditionally non-exercising situation include, but are not limited to the following: 1) allowing low-impact office and/or therapy exercise; 2) improving blood circulation; and 3) enhancing weight control. These benefits would greatly help people who either need or want to exercise more and those who are in medical need of exercise through either a therapeutic or rehabilitation program that requires them to exercise during the day, regardless if they are sitting or not.\nFurthermore, the mode of motion that a person moves his or her appendages to exercise is important. If one is sitting and performing a stepping motion with his or her feet, one is only exercising a certain subset of muscles. Likewise a different but limited subset of muscles is used for a cycling (circular or eliptical) or sliding (back and forth) motions. Additionally, being capable of alternating between different modes of exercising motion is important not only to exercise more and different muscles and muscle groups, but also to prevent an exercising person from becoming bored and uninterested in continuing to exercise.\nPrior art exercise devices for use while sitting at a desk or similar locale commonly employ a stepping motion of a person's feet or at best one mode of motion for the feet, such as stepping, cycling, or sliding. However, the prior art does not appear to combine multiple modes of motion to permit one to exercise in their choice of motional mode.\nAn example of prior art that involves a stepping motion while sitting in a chair includes U.S. Pat. No. 6,709,368 issued to Chue, Mar. 23, 2004, for Foot Exercise Device. This patent describes a device that has two steps (i.e. pedals or levers), as the exerciser is of the stair or step variety. First one foot presses down on one step, as the other step rises. Then, the other foot pushes down on the other step, and the stepping motion is reversed. Variable resistance in the form of minilevers is provided.\nAnother example of prior art that involves a stepping motion while sitting in a chair includes U.S. Pat. No. 6,042,521 issued to De Giorgis, Mar. 28, 2000, for Exercising Means. This patent describes an exercise device to enable a person to exercise his or her feet while seated. More particularly, the exercising device enables a sitting person to simulate a walking action. The exercising device has a base, a pivot beam, and two footrests. Rear parts of the footrests can pivot with respect to the base by way of ball and socket fittings, and forward parts of the footrests can pivot with respect to the pivot beam by way of additional ball and socket fittings. The pivot beam can pivot with respect to the base by way of a further ball and socket fitting. When the exercising device is in use the person can place his or her feet on the footrests and tilt both the feet and footrests from side to side. The person can also cause forward parts of the footrests to move up and down alternately by way of the pivot beam by alternately pushing his or her feet down against the footrests.\nAn example of prior art that involves a cycling motion while sitting in a chair includes United States Patent Application Pub. No. 2002/0107114 issued to Byrd, Jr., Aug. 8, 2002, for Rotational Exercise System, Device and Method. This patent application publication describes an exercise device that includes a frame, a cranking assembly mounted to the frame and an extended user fixture. In one implementation of the device, a user sits in a chair and places his or her feet on pedals and moves the pedals in a cycling motion to turn a flywheel.\nAnother example of prior art that involves a cycling motion while sitting in a chair includes the Chattanooga Deluxe Exerciser, offered for sale by Promed Products of Atlanta, Ga., USA. This portable device requires a sitting user to place his or her feet on the footrests and then pedal in a cycling motion. The device has electronic controls and an LED display that shows you information such as speed, distance, total mileage, total time “ridden”, and the amount of calories burned.\nAn example of prior art that involves a sliding motion while sitting in a chair includes U.S. Pat. No. 5,833,575 issued to Holslag, Nov. 10, 1998, for Portable Exercise Apparatus having Chair Mountable Support Base and Variable Resistance Exercise Arms. This patent describes a portable exercise apparatus that includes a support frame having a base mountable under a chair to hold the base in a stationary position with a user seated on the chair and an upright standard mounted upon the base, a shaft mounted to an upper end of the upright standard and having opposite ends extending from opposite sides thereof, a pair of arms disposed on opposite sides of the upright standard with each arm at one end mounted to one end of the shaft for rotatably mounting the arm to the upright standard, a pair of pedals each mounted to the other end of each of the arms for engagement by a user to create the force necessary to rotate the arms relative to the upright standard, and a resistance generating and adjusting mechanism disposed on the ends of the shaft at the opposite sides of the upright standard and engaged with the one ends of the arms and being operable to generate and selectively adjust a level of resistance to rotation of the pair of arms relative to the upright standard in response to rotation of the arms.\nAn example of prior art that attempts to combine multiple modes of motion while sitting in a chair includes the Zen-Pro Foot and Leg Exerciser, offered for sale by Zen-Pro Ltd of Shillingford, United Kingdom. This device requires the user to place his or her feet on the footrests, then one may rotate the feet on the rests and/or rock the device side-to-side. The device is not stable as the base is curved in order to accommodate the rocking motion. Furthermore, the device's foot holders are connected directly to the base, through a ball-in-socket connector.\nThe prior arts do not appear to successfully combine at least the three previously discussed modes of exercising motion: sliding, cycling, and stepping. Thus, there remains a need for a compact low impact therapeutic exercise device designed for increasing circulation in a sitting user having multiple modes of motion."}
{"text": "Embodiments of the invention relate generally to information retrieval in a network and, more particularly, to hosting and distributing content on a content delivery network such as the Internet.\nThe World Wide Web is the Internet's content retrieval system. In the Web environment, client systems effect transactions to Web servers using the Hypertext Transfer Protocol (HTTP), which provides clients with access to files (e.g., text, graphics, images, sound, video, etc.) using a standard page description language, for example, Hypertext Markup Language (HTML). A network path to a server is identified by a so-called Uniform Resource Locator (URL) having a special syntax for defining a network connection. Use of a Web browser at a client (end user) system involves specification of a link via the URL. In response, the client system makes a request to the server (sometimes referred to as a “Web site”) identified in the link and, in return, receives content from the Web server. The Web server can return different content data object types, such as .gif and jpeg files (graphics), .mpeg files (video), .wav files (audio) and the like.\nAs the Web server content provided by World Wide Web has continued to grow over time, so too has the number of users demanding access to such content. Unfortunately, the ever-increasing number of end users requesting Web content from Web sites has resulted in serious bandwidth and latency issues, which manifest themselves in delay to the end user.\nTo address these problems, many networking product and service providers have developed solutions that distribute Web site content across the network in some manner. One class of solutions involves replicating Web servers at multiple locations and directing traffic (by modifying the URL and forwarding, or using HTTP re-direct) to the “best” server based on a predefined selection policy, e.g., load balancing, network topology. Another class of solutions distributes content strategically and/or geographically, and often uses some type of centralized or hierarchical Domain Name System (DNS)-based site selection. The distributed sites include servers that perform reverse proxy with (or without) caching. One such technique routes traffic to a content distribution site nearest the requestor by modifying URLs in the top-level Web page. Other DNS techniques use a round robin traffic distribution to distribute load to the content sites, but do not take into account the location of the requester relative to those content sites."}
{"text": "The need for information cards that contain electronics such as integrated circuits (IC) has been increasing as of late. The IC can provide functions, such as for retrieving, processing, transmitting and storing information. Such cards are generally referred to as “smart cards” which have potential application in a wide range of fields, including commerce, transportation, communication, identification, and security.\nIt is known to incorporate batteries into cards to provide a power source for the electronic circuitry within. Having a battery incorporated into the IC card allows for the storage of greater amounts of data, and further allows for improved processing capabilities.\nBecause of the difficulty of providing the energy required for most applications within the space constraints of a typical credit card, it is desirable to use battery types which provide the highest energy density. Lithium batteries have been introduced into the market because of their high energy densities. Lithium is atomic number three on the periodic table of elements, having the lightest atomic weight and highest energy density of any solid material. As a result, lithium is a preferred material for batteries, having very high energy density. Lithium batteries are also desirable because they have a high unit cell voltage of up to approximately 4.2 V, as compared to approximately 1.5 V for both Ni—Cd and Ni-MH cells. Lithium batteries can be either lithium ion batteries or lithium metal batteries. Lithium ion batteries intercalate lithium ions in a host material, such as graphite, to form the anode. On the other hand, lithium metal batteries use metallic lithium or lithium alloys for the anode.\nLithium batteries having solid polymer electrolytes represent an evolving alternative to lithium batteries having liquid electrolytes. Solid polymer electrodes are generally gel type electrolytes which trap solvent and salt in pores of the polymer to provide a medium for ionic conduction. Typical polymer electrolytes comprise polyethylene oxide (PEO), polyether based polymers and other polymers which are configured as gels, such as polyacrylonitrile (PAN), polymethylmethacrylate (PMMA) and polyvinylidine fluoride (PVDF). The polymer electrolyte generally functions as a separator, being interposed between the cathode and anode films of the battery.\nLithium metal polymer (LMP) rechargeable batteries offer improved performance as compared to Li ion batteries, particularly higher capacity. LMP batteries result from the lamination/assembly of three types of main thin films: a film of positive electrode comprising a mixture of a polymer and an electrochemically active material such as lithium vanadium oxide, an electrolyte film separator made of a polymer and a lithium salt, and a negative electrode film comprising metallic lithium or a lithium alloy.\nUnfortunately, conventional batteries including lithium metal or lithium ion batteries lack the stability to be subjected to credit card processing, which can include a temperature of 125C to 140C, a pressure of 200 to 250 psi, and a dwell time of 5 to 15 minutes. Conventional batteries subjected to the above described card lamination conditions will generally either catastrophically fail or experience a large shift in operating parameters which would render them of no practical use. As a result, although batteries have been disposed inside credit cards, a low temperature glue has been used to bind the battery to the card."}
{"text": "This invention relates to a device for once-only use, preferably for application to hygiene.\nIn the field of the aforementioned devices and in particular those used for administering vaginal washes or enemas, devices for once-only use are known comprising a deformable container, to the mouth of which a cannula can be sealedly fixed for insertion into the patient's body.\nKnown devices for once-only use are sold with the container already filled with the solution ready for use and with their mouth sealed by a stopper which is removable, preferably manually. For hygienic reasons the cannula is packaged separately in a sealed packet, but sold together with the container.\nThe user (after necessarily shaking the container to revive the solution) removes the stopper from the mouth, withdraws the cannula from its package and fits it to the container mouth to obtain a device ready for use.\nThe container of such devices is made of plastics material, as it has necessarily to be impermeable, soft and flexible.\nExperience has, however, shown that plastics material is not the best material for preserving certain substances which, being ready for use, are particularly active and hence aggressive and/or instable, and consequently of limited life.\nTo overcome this drawback, devices have been considered comprising two containers of different capacity, plus the already known cannula. One container is filled with some of the components of the solution and the other is filled with the remaining components. The division is obviously made so as to reduce the chemical activity of the components to a minimum, so achieving low aggressivity, substantial chemical stability and a long life. Such devices can therefore be constructed of plastics material at low cost.\nForming the device as two containers can also be dictated by other requirements known to the expert but not mentioned herein for brevity.\nTo form the solution ready for use the user has to pour the contents of the smaller container into the larger container in addition to carrying out all the operations already mentioned in relation to the preceding case, i.e., removing the stopper from two containers and fitting the cannula. There is also an obvious danger in pouring substances which usually have a high capacity for impregnation and hence a high staining power.\nIn addition, this type of device does not allow optimum mixing of the contents of the two containers if the contents are poorly miscible with each other. In such a case, the solution ready for use can only be achieved in an open container closable by the user's fingers during shaking, with the possibility both of contaminating the product and producing accidental splashing.\nHence, in the current state of the art, the devices in use may comprise one or two containers. The choice is obviously dictated by the type of solution to be sold."}
{"text": "In 2010, the WHO estimated that there were 285 million globally with visual impairment, 39 million of who were “blind” and 246 million who had “low vision.” A traditional white cane is a common tool used by blind or low vision people. James Biggs, a photographer blinded following an accident, invented the white cane in 1921; he found it prudent to paint his walking stick white as a symbol of his impairment for the heavy traffic surrounding his domicile. This movement was later fostered by Guilly d'Herbemont of France in 1931 for the French veterans of WWI. Individuals and companies have attempted to further the advancement of the white cane or navigating tools. However, at the end of the day prosthetics departments offer amputees complex $70,000 joint components and our visual occupational therapists offer our legally blind a simple, modified, low profile, white cane tip. Simply observing a visually impaired individual ambulate with a conventional white cane over several minutes on a Manhattan street provides enough encouragement to redefine their assistive device."}
{"text": "This invention relates to a heat-gain structure. It is illustrated as a high efficiency cold frame, although its usefulness is not limited thereto.\nThe need for providing a structure which will increase the growing season for plants has been long felt. One structure for prolonging the growing season is the cold frame or hot bed.\nA cold frame consists of a frame and a cover glazed with a light-transmitting material. The cold frame provides a heat gain during daylight, even on most cloudy days, and provides some retention of heat during cold and sunless periods. Cold frames, however, have a number of well-known drawbacks.\nOn sunny days, cold frames tend to overheat, even in very cold weather. To overcome this problem, temperature responsive openers have been used, as shown, for example, in Anderson, U.S. Pat. No. 3,129,530, and in Caldwell, U.S. Pat. No. 4,005,547. This expedient has not been entirely successful, because opening the cold frame sufficiently to let the heat out has sometimes permitted cold wind to freeze plants, especially near the edges of the cold frame.\nWind causes other problems with cold frames. When the glazing is laid back to open the frame, it is subject to being blown away and broken by the March or October wind. If the glazing is hinged to the frame, it may well carry the entire frame with it or deform the hinges or frame. When the glazing is partially opened, either manually or by an automatic opener, it is even more subject to being caught by the wind than when it is laid back on the ground.\nCold frames are also not particularly efficient heat retainers. On cold nights, the plants within them are frequently damaged. To prevent heat loss, a blanket of insulation is sometimes placed over the cold frame, and a water-resistant sheet is placed over the blanket. The blanket and sheet must then be secured against wind, because if even a corner of the blanket blows loose, its insulation value is largely lost.\nThe time and physical effort required to manipulate cold frames during harsh weather, as well as their limited effectiveness, has limited their use.\nHot beds are essentially cold frames with a supplemental source of heat, either from an outside energy source or from decaying matter. Although they reduce the problem of freezing, they are inefficient users of energy and suffer all of the other problems associated with cold frames.\nGreenhouses are much larger, hence much more expensive, than cold frames. An unheated greenhouse loses heat rapidly. A heated greenhouse requires great amounts of energy to provide supplemental heating at night, during cloudy or extremely cold weather. The cost of energy for heating greenhouses has become a serious problem. Various methods for insulating greenhouses have been proposed, but they have generally been so expensive or cumbersome that they have not been widely adopted."}
{"text": "1. Field of the Invention\nThis invention generally relates to the use of a network device/appliance to store its own newest driver, or other desired device driver desired by a user and to assist other computers or devices to update their drivers for using the device/appliance. This invention is further related to configuring a service center or resource administration station to download the newest or desired driver to the network device/appliance remotely.\n2. Discussion of the Background\nThe cross-referenced applications and patents listed above show various methods to communicate with networked devices and appliances. In particular, a service center can establish communication via either direct communication (e.g., a connectionless (UDP) socket or connection-based (TCP) socket using an information protocol (e.g., an HTTP or FTP), or indirect communication (e.g., a store-and-forward system such as electronic mail (e-mail))).\nBusiness office appliances such as copiers, facsimiles, multi-function image processors, printers and image storage systems have been connected to networks recently. However, in order to use the business office appliances, the device drivers of the appliances need to be installed at client computers. As newer drivers become available, it is necessary to update the driver of each client computer to take full advantage of the functions added by the newer driver."}
{"text": "The invention herein described relates generally to friction test machines and methods and, more particularly, to a friction testing machine and method for measuring the coefficient of friction between a rubber specimen or a tread element and different friction surfaces at different speeds, contact pressures and orientations.\nWhen developing new compounds for tires, a prior practice was to build a tire and subject the tire to various traction tests. Because of the cost and time needed to build test tires for the purpose of optimizing traction properties, small samples of tread compounds have instead been used to determine the traction characteristics of the compounds. By using small samples for testing, a large number of different tread compounds can be made in small batches for screening to determine which compound has the best properties. The small samples could also be tested with different friction test surfaces representing different road surfaces and different conditions, whereby the traction properties of the compounds may be observed to determine which compound has the best traction on a specific road surface under specific conditions.\nIn U.S. Pat. No. 5,113,688 there is disclosed an apparatus and method for predicting tire traction characteristics of tread compounds using small test samples. The apparatus causes a rotating relationship between the sample and friction surface. The torque between the sample and friction surface is measured and a torque versus time curve is established for the sample. Such apparatus, however, does not appear to be easily movable as may be desired for testing at different locations, nor is such apparatus suitable for evaluating the performance of different tread patterns. It would be advantageous to have a portable friction testing apparatus that not only can measure the coefficient of friction of plain tread elements, but also can measure the coefficient of friction of tread blocks, shoulders, etc., of existing tires. A further desirable advantage would be an apparatus that permits one to view the interaction between the friction surface and the specimen in the same way as the present invention.\nThe present invention provides a friction testing machine and method for measuring friction characteristics between a test sample and a friction surface. The machine and method are particularly suited for measuring the coefficient of friction between a rubber specimen or a tread element and different friction surfaces at different sliding velocities, contact pressures and orientations. A preferred embodiment of machine is self-contained and portable, configured for easy and quick changing of the friction surface, and provides for rotating the test sample about an axis normal to the sliding surface and the direction of movement of the sample relative to the friction surface.\nIn general, a friction test machine according to the invention comprises a sample holder configured to hold a sample in frictional engagement with a friction surface, and a motion device for effecting relative movement between the friction surface and sample holder in a first direction. Also provided is a variable weight loading device on the sample holder for loading the sample holder such that a selected load can be applied to the sample normal to the friction surface, and a force measurement device for obtaining a measurement indicative of the frictional force resisting such relative movement between the sample holder and the friction surface effected by the motion device.\nPreferably, there is provided a second motion device for effecting relative movement between the sample holder and friction surface in another direction different than the first direction, such as rotation about an axis perpendicular to the friction surface and/or translation in a direction transverse to the primary direction of movement.\nIn one embodiment of friction test machine according to the invention, there is provided a base on which the friction surface is supported, and a carriage mounted on the base and guided for movement in a first direction parallel to the friction surface. The motion device is connected between the carriage and base for moving the carriage in the first direction, and the sample holder is mounted to the carriage for movement therewith and configured to hold the sample in frictional engagement with the friction surface during movement of the carriage in the first direction. In another embodiment of the friction test machine, the motion device causes relative movement between the carriage and the friction surface in a second direction opposite to the first direction so that the force measurement device can obtain measurements indicative of the frictional force in forward and reverse directions.\nPreferably, the friction surface is provided on a friction element, and the friction element is removably supported on the base, whereby the friction element can be interchanged with other friction elements for providing respective different friction surfaces. In a preferred embodiment, the base includes a recess for removably holding the friction element, and the base is supported on a cart for easy transport of the machine from place to place.\nIn another embodiment of friction testing machine according to the invention, there is provided a carriage fixed to the base, and the sample holder is mounted in the carriage for vertical movement perpendicular to the friction surface. A table is mounted to the base for rotation about a vertical axis, and the friction surface extends annularly around a radially outer peripheral portion of the table, whereby the friction surface is moved while the sample remains stationary.\nIn still another embodiment of friction testing machine according to the invention, there is provided a base adapted to rest on a surface against which a sample is to be tested, such as a road surface. A carriage is mounted on the base and guided for movement in a first direction parallel to the friction surface. A motion device is connected between the carriage and base for moving the carriage in the first direction, and a sample holder is mounted to the carriage for movement therewith and configured to hold the sample in frictional engagement with the friction surface during movement of the carriage in the first direction. Preferably, the base is provided with an aperture through which the sample holder extends for positioning the sample against the friction surface disposed beneath the base.\nIn a preferred embodiment of the invention, the variable weight loading device includes a weight platform on which removable weights can be stacked and removed to selectively vary the load applied to the sample. Alternatively, the variable weight loading device comprises a fluid cylinder that is controlled by the processor to vary the load applied to the sample. It also is preferred to provide the friction surface on a friction element that is easily and quickly interchangeable with other friction elements for providing respective different friction surfaces.\nIn a preferred embodiment of the invention, a temperature chamber encloses the friction surface and the sample whereby temperature can be controlled to simulate different temperature conditions. Additionally or alternatively, the machine includes a water source and a pump which transfers water from the source onto the friction surface whereby the measurement device can obtain measurements indicative of the frictional force in wet conditions.\nIn a preferred embodiment of the invention, a processor controls the motion device, controls the variable weight loading device, controls the temperature chamber, and/or records the measurements obtained by the force measurement device. The processor further analyzes the measurements and compiles the results in the form of tables and/or graphs and an output device is provided to allow a readout of the compiled results. Preferably, the processor provides a library of pre-defined tests and a selection device for selecting a particular test.\nThe present invention accordingly provides a method of predicting tire traction characteristics of a tread component. The method comprises the steps of preparing a sample of the tread component; attaching the sample to the sample holder, setting the variable weight device to apply a certain load to the sample, and using the force measurement device for obtaining a measurement indicative of the frictional force resisting such relative movement between the sample holder and the friction surface effected by the motion device. During this method, the friction surface may be replaced with another different friction surface to predict tire traction characteristics with respect to different road surfaces. Also, water may be transferred onto the friction surface to gather measurements indicative of the frictional force in wet conditions. The temperature surrounding the friction surface and the sample may be controlled to simulate different temperature conditions, such as winter conditions and summer conditions.\nThe foregoing and other features of the invention are hereinafter fully described and particularly pointed out in the claims, the following description and the annexed drawings setting forth in detail one or more illustrative embodiments of the invention, such being indicative, however, of but one or a few of the various ways in which the principles of the invention may be employed."}
{"text": "Higher Order Ambisonics denoted HOA is a mathematical description of a two- or three-dimensional sound field. The sound field may be captured by a microphone array, designed from synthetic sound sources, or it is a combination of both. HOA can be used as a transport format for two- or three-dimensional surround sound. In contrast to loudspeaker-based surround sound representations, an advantage of HOA is the reproduction of the sound field on different loudspeaker arrangements. Therefore, HOA is suited for a universal audio format.\nThe spatial resolution of HOA is determined by the HOA order. This order defines the number of HOA signals that are describing the sound field. There are two representations for HOA, which are called the spatial domain and the coefficient domain, respectively. In most cases HOA is originally represented in the coefficient domain, and such representation can be converted to the spatial domain by a matrix multiplication (or transform) as described in EP 2469742 A2. The spatial domain consists of the same number of signals as the coefficient domain. However, in spatial domain each signal is related to a direction, where the directions are uniformly distributed on the unit sphere. This facilitates analysing of the spatial distribution of the HOA representation. Coefficient domain representations as well as spatial domain representations are time domain representations."}
{"text": "While just about every computer user owns their own printer and is capable of producing high quality documents, the ability to produce such documents in high volume and with special finishing features, such as binding, is still within the purview of the commercial print shops and corporate copy departments. High volume, finished production of documents is typically referred to as production printing. A production printer is a printing device capable of rapid production of large volumes of documents. Typically these printers have high paper handling capacity, the ability to draw on multiple media types from multiple sources and the ability to automatically finish a document such as by adding a binding. Despite the automation provided by the production printer and the proliferation of computer technology, especially in the area of desktop publishing, production printing is still a complicated and often manual process.\nIn a typical print shop, customers bring in original documents which they want turned into a finished product such as a bound booklet, a tri-fold brochure or a tabbed three ring bound notebook. In addition, they typically need a large volume of the finished product, for example, one thousand brochures. The combination of the original documents plus the instructions for producing the finished product is called a “job”. The documents can be brought in either in hard copy or electronic form, such as on floppy disk, compact disc or tape or can be transmitted to the print shop over a network such as the Internet.\nAfter handing over the documents to the clerk, the customer relays his instructions for preparing the finished product. The clerk will note these instructions on a “ticket” or “job ticket”. The job ticket is typically a piece of paper with all of the instructions written on it for producing the finished product. As mentioned above, this is known as job. The job will then be handed to an operator, who runs the production printer, to produce the finished output. The operator's job is to prepare the document for production, load the appropriate materials, such as paper stock and binding materials, into the production printer and ensure that the finished output is correct.\nWhile the job of the operator seems simple, there are many issues which quickly complicate it. Often, the documents provided by a customer are not ready to be run on the production printer. Some documents provided by a customer are merely raw manuscripts requiring basic formatting, such as margins, typography, etc. Other documents may be formatted but such formatting might not take into account the requested binding. For example, the text of the document is too close to the margin, therefore, when the finished product is bound, some of the text will be obscured. Some documents, such as books, require special care so that, for example, the first page of every chapter appears on the front of a page, also known as imposition. Other forms of imposition include booklet/pamphlet imposition or n-up imposition. Or the customer may bring in multiple documents and ask that these “chapters” be assembled into a book, with a cover and binding.\nOther issues which complicate the production printing job are determining and loading the correct media into the production printer. Often, jobs will require many different paper types, such as different stock weights or different colors. In addition, some jobs require the insertion of tab stock at specific points within the document. Still other jobs may require the adding of a Bates number or other annotation to the document.\nWith such a complicated production process to produce finished output, errors are bound to occur, such as loading the wrong paper stock in the printer or setting a margin too close to a binding. Production printers run at very high speeds, often producing output greater than 1 page per second therefore, errors in the finished output may not be caught before a significant amount of time and resources have been wasted.\nAccordingly, there is a need for an efficient system and method for managing the production printing workflow."}
{"text": "During recent years, computer gaming has gained increasing popularity, and today thousands of players are playing games all around the world. Predictably, interactive computer gaming has blurred the line between games and other entertainment or communication media, and the avenues that are being explored in the development of gaming might well break new ground for interactive Internet applications in all areas of business relations and social life.\nDue to their dynamic nature and specific appeal to certain audiences, computer games and especially games played on Internet enabled platforms provide the ideal vehicle for Internet advertising. Not only is it possible for an advertiser to target directly a specific group of customers, but also advertisements can be directly incorporated into the computer games, enabling something akin to the well-known concept of product placement. In such a system, a game is running on a system that is connected to the Internet. Throughout the playing of the game, the game software contacts a set of ad servers to retrieve advertising content to present to the user during game play. The data and advertising content is transferred using HTTP protocols over the Internet. Optionally, the advertising content is cached locally to the game, and throughout the playing of the game the game software contacts a set of ad servers to retrieve a list of advertising content to present to the user during game play.\nFor product placement advertisements, there is a considerable development effort involved in their implementation and a considerable business effort in sales, customer service, and reporting. As such and because of the constant communication, each party involved in the advertisement design, placement, and integration is provided an opportunity to present input opinions and decisions. With all of this communication, advertisements are approved by each party as part of their involvement. From the game developer, to the game producer, to the advertiser, to the game console manufacturer, to third party IP holders each party sees and approves of the content through their business dealings. Clearly, for one or two advertisements built into a video game, this is manageable.\nWith the advent of dynamic in-game advertising, there is significantly more effort in garnering approval for advertising content and yet, companies involved in in-game advertising continue to pursue manual approvals in order to ensure that all parties approve of advertising content for a video game. Such a process is slow, costly, and prone to human error, where the error is often uncorrelated to the parties involved—one party makes an error that adversely affects another party. This is very undesirable.\nOf course, one of the great problems that in-game advertising is to face is that for a single product placement, relatively senior people within an organization involve themselves in the process. For 1000 advertisements, the hands on executive is no longer feasible. Unfortunately, with more advertisements and fewer senior level people involved, a recipe for disaster results.\nIt would be advantageous to provide a method and system that overcomes these and other drawbacks of the prior art.\nIt is an object of the invention to provide an enhanced content approval process for use with in-game advertising.\nIn accordance with the invention there is provided a method comprising: providing a request for approval of advertising content for being impressed upon video game players within a known advertising spot within a known video game; determining based on the advertising spot and the video game at least a party for approving of the advertising content; and automatically transmitting a message to the at least a party requesting their approval of the advertising content for being impressed upon video game players within the known advertising spot within the known video game.\nIn accordance with the invention there is provided a method comprising: receiving secured advertising content from a remote server, the advertising content secured against tampering; determining from the secured advertising content an approval status thereof; determining based on the secured advertising content an absence of tampering thereof; and, when the secured advertising content is both approved and absent tampering, storing of the advertising content in an accessible form within a memory storage for retrieval during video game execution and for impressing upon a gamer of the video game.\nIn accordance with another aspect of the invention there is provided a system comprising: a suitably programmed processor for: providing a request for approval of advertising content for being impressed upon video game players within a known advertising spot within a known video game; determining based on the advertising spot and the video game at least a party for approving of the advertising content; and automatically transmitting a message to the at least a party requesting their approval of the advertising content for being impressed upon video game players within the known advertising spot within the known video game.\nIn accordance with another embodiment of the invention there is provided a system comprising: a suitably programmed processor for: wherein secured advertising content is received at a video game system; wherein software in execution on the video game system determines from the secured advertising content an approval status thereof; wherein software in execution on the video game system determines based on the secured advertising content an absence of tampering thereof; and, wherein storing of the advertising content comprises storing of the advertising content in an accessible form within a memory storage of the video game system for retrieval during video game execution and for impressing upon a gamer of the video game.\nIn accordance with another aspect of the invention there is provided a storage medium having data stored therein, the data for when executed resulting in performance of: providing a request for approval of advertising content for being impressed upon video game players within a known advertising spot within a known video game; determining based on the advertising spot and the video game at least a party for approving of the advertising content; and automatically transmitting a message to the at least a party requesting their approval of the advertising content for being impressed upon video game players within the known advertising spot within the known video game.\nIn accordance with another embodiment of the invention there is provided a storage medium having data stored therein, the data for when executed resulting in performance of: a suitably programmed processor for: wherein secured advertising content is received at a video game system; wherein software in execution on the video game system determines from the secured advertising content an approval status thereof; wherein software in execution on the video game system determines based on the secured advertising content an absence of tampering thereof; and, wherein storing of the advertising content comprises storing of the advertising content in an accessible form within a memory storage of the video game system for retrieval during video game execution and for impressing upon a gamer of the video game."}
{"text": "1. Technical Field of the Invention\nThe present invention relates to a method and apparatus of the positioning of a fluid actuator and particularly to improved method and apparatus for rapidly and accurately setting a fluid actuator by supplying pressurized fluid to the right and left actuating chambers thereof through an on-off control valve.\n2. Description of the Prior Art\nHydraulic and pneumatic actuators for driving or positioning various types of operating devices are well-known and currently utilized as fluid actuators in various fields of industry. Although these fluid actuators have been utilized mainly in industrial machines, they are recently being used to control robots. There is accordingly required a fluid actuator which can rapidly and accurately be positioned under severer conditions.\nIn general, the fluid actuator comprises an on-off control valve through which a pressurized fluid is supplied to the right and left actuating chamber of the actuator, the desired position of the actuator being depending on the amount of the supplied fluid. The relatively rough construction of the prior art fluid actuators are disadvantageous in that the positioning operation cannot rapidly be attained within a predetermined range of movement with an accuracy of 100 .mu.m or less only by operating the control valve.\nTo overcome such a disadvantage in the prior art, some improved fluid actuators have been proposed. For example, Japanese Laid-Open Patent Sho No. 58-174701 discloses a fluid actuator which includes an electromagnetic valve having an expectable delay in response and hydraulic cylinder controlled in accordance with a predetermined result of operation without hunting which would easily be created in such a fluid actuator. However, this proposal also cannot accomplish any high-accuracy positioning."}
{"text": "The invention relates to a method and arrangement for determining and indicating the position and depth of an unevenness of the inner surface of a pipe.\nSuch method is for example disclosed in German patent application No. P 30 03415, corresponding to U.S. patent application Ser. No. 230,050, filed Jan. 29, 1981, now U.S. Pat. No. 4,407,071, which process has the drawback that the mechanical scanning system comes into contact with the pipe wall. In view of the frictional forces that are encountered, errors in the measurements and even jamming of the feelers in imperfections or erosions of the interior surface of the pipe, can result. Due to the mechanical scanning, the period between two measurements is relatively long since the feeler finger must be lifted off the interior pipe wall, shifted in position and then again brought into contact with the pipe wall."}
{"text": "The present invention is directed toward a movable roof structure and more particularly toward a telescoping movable roof structure for a patio, swimming pool, building or the like.\nNumerous types of facilities such as swimming pools, tennis courts and patios are frequently used in the open air during fair weather. During inclement weather, however, it would be desirable to provide such facilities with a cover or roof. In order to fulfill both desires, roofs have been proposed in the past which are collapsible or movable so that a closed top can be provided if desired or the same can be opened. Such movable roofs are also sometimes desirable with various other types of buildings.\nThe proposed prior art movable roof structures known to Applicant are not believed to have been satisfactory. They have either been too complex and, therefore, too difficult and expensive to construct or have been capable of opening only a small portion of the roof.\nFor example, while the structures shown in U.S. Pat. Nos. 3,566,555; 3,589,084 and 4,073,098 are not very complex, they provide for the movement of only a single panel member or section and accordingly only a portion of the roof can be opened. The structure shown in U.S. Pat. No. 4,283,889 is capable of opening the entire roof but only because the roof is relatively small as only one roof panel is also moved in this structure. U.S. Pat. Nos. 1,006,374 and 2,094,801 show movable roof structures comprised of a plurality of panels which are adapted to telescope with respect to each other so that a larger portion of the roof can be opened. This is a similar concept to Applicant's invention. However, the structure shown in these two latter patents can be used only on a particular predetermined structure. They can be adapted to other structures only by making changes to the roof system itself or the supporting structure or both. This, of course, increases the cost of the roofing system since each must be specially designed."}
{"text": "Capturing of acoustic information is known to be performed from a static location. In addition, the use of cameras on aerial vehicles to collect information is known. But, significant technical challenges exist in sensing acoustic information from aerial vehicles. In particular, due to rigorous size, weight, and power requirements and corresponding technical challenges necessary to integrate acoustic sensing payload hardware, no related art system can reliably capture voice communication from an unmanned aerial vehicle (UAV). Therefore, a need exists for an acoustic sensing payload that is capable of collecting acoustic data from a UAV."}
{"text": "Golf is enjoyed by a wide variety of players—players of different genders and dramatically different ages and/or skill levels. Golf is somewhat unique in the sporting world in that such diverse collections of players can play together in golf events, even in direct competition with one another (e.g., using handicapped scoring, different tee boxes, in team formats, etc.), and still enjoy the golf outing or competition. These factors, together with the increased availability of golf programming on television (e.g., golf tournaments, golf news, golf history, and/or other golf programming) and the rise of well known golf superstars, at least in part, have increased golf's popularity in recent years, both in the United States and across the world.\nGolfers at all skill levels seek to improve their performance, lower their golf scores, and reach that next performance “level.” Manufacturers of all types of golf equipment have responded to these demands, and in recent years, the industry has witnessed dramatic changes and improvements in golf equipment. For example, a wide range of different golf ball models now are available, with balls designed to complement specific swing speeds and/or other player characteristics or preferences, e.g., with some balls designed to fly farther and/or straighter; some designed to provide higher or flatter trajectories; some designed to provide more spin, control, and/or feel (particularly around the greens); some designed for faster or slower swing speeds; etc. A host of swing and/or teaching aids also are available on the market that promise to help lower one's golf scores.\nBeing the sole instrument that sets a golf ball in motion during play, golf clubs also have been the subject of much technological research and advancement in recent years. For example, the market has seen dramatic changes and improvements in putter designs, golf club head designs, shafts, and grips in recent years. Additionally, other technological advancements have been made in an effort to better match the various elements and/or characteristics of the golf club and characteristics of a golf ball to a particular user's swing features or characteristics (e.g., club fitting technology, ball launch angle measurement technology, ball spin rates, etc.). Further technological advancement in golf club design has also involved the incorporation of various types of monitoring devices or sensors in the golf club. Many such designs, however, have been cumbersome and unreliable. In addition, further processing of the data recorded by the sensors has been limited or not performed in a suitable manner to be most useful to golfers.\nWhile the industry has witnessed dramatic changes and improvements to golf equipment in recent years, there is room in the art for further advances in golf club technology. Thus, while golf equipment according to the prior art provide a number of advantageous features, they nevertheless have certain limitations. The present invention seeks to overcome certain of these limitations and other drawbacks of the prior art, and to provide new features not heretofore available."}
{"text": "Cameras are widely used by average consumers, hobbyists, and professionals. While film-based photography has been popular for decades, digital cameras are quickly becoming the principal photography device for most households. Indeed, many electronic devices are being integrated with digital camera functionality, and traditional film cameras are being fitted with electronic features as well. For example, mobile telephones which include a digital camera are becoming increasingly common.\nA well known photographic method involves utilizing a different depth of field (“DOF”) to produce an image with the subject slightly out of focus, which is sometimes called “soft-focus” or portrait-style. Depth of field is the distance between the farthest and nearest points which are in focus. This can also be identified as the zone of acceptable sharpness in front of and behind the subject, to which the lens is focused on. While “acceptably sharp” is a somewhat subjective standard, it has been defined mathematically by those skilled in the art. In general, the suggested photographic procedure is to place ⅓ of the DOF in front of the subject and ⅔ behind the subject.\nThe DOF is determined from several variables and is a function of the size of the image capturing device (i.e. film or photoelectric sensor), aperture, focal length, and distance to the subject. In portrait photography, it is common to utilize a shallow DOF in order to place the subject in sharper contrast to the surroundings.\nFilm-based cameras have utilized various methods such as special lenses, filters, or even gels to achieve a shallower DOF. However, since DOF is inversely proportional to the size of the image capturing device, film-based cameras have traditionally had much greater control over DOF than digital cameras, due at least in part to the relative large size of film in comparison to photoelectric sensors used in digital cameras. Thus, the relatively small size of the image capturing device used in digital cameras has traditionally resulted in a large DOF. Therefore, there is a need for a method and apparatus to allow for a shallower DOF, such that a user can employ various photographic techniques such as portrait photography."}
{"text": "Object-oriented technology (OOT), which generally includes object-oriented analysis (OOA), object-oriented design (OOD), and object-oriented programming (OOP), is earning its place as one of the most important new technologies in software development. OOT has already begun to prove its ability to create significant increases in programmer productivity and in program maintainability. By engendering an environment in which data and the procedures that operate on the data are combined into packages called objects, and by adopting a rule that demands that objects communicate with one another only through well-defined messaging paths, OOT removes much of the complexity of traditional, procedure-oriented programming.\nThe following paragraphs present a brief overview of some of the more important aspects of OOT. More detailed discussions of OOT are available in many publicly available documents, including Object Oriented Design With Applications by Grady Booch (Benjamin/Cummings Publishing Company, 1991) and Object-Oriented Requirements Analysis and Logical Design by Donald G. Firesmith (John Wiley & Sons, Inc., 1993). The basic component of OOT is the object. An object includes, and is characterized by, a set of data (also called attributes) and a set of operations (called methods) that can operate on the data. Generally, an object's data may change only through the operation of the object's methods.\nA method in an object is invoked by passing a message to the object (this process is called message passing). The message specifies a method name and an argument list. When the object receives the message, code associated with the named method is executed with the formal parameters of the method bound to them corresponding values in the argument list. Methods and message passing in OOT are analogous to procedures and procedure calls in procedure-oriented software environments. However, while procedures operate to modify and return passed parameters, methods operate to modify the internal state of the associated objects (by modifying the data contained therein). The combination of data and methods in objects is called encapsulation. Perhaps the greatest single benefit of encapsulation is the fact that the state of any object can only be changed by well-defined methods associated with that object. When the behavior of an object is confined to such well-defined locations and interfaces, changes (that is, code modifications) in the object will have minimal impact on the other objects and elements in the system. A second \"fringe benefit\" of good encapsulation in object-oriented design and programming is that the resulting code is more modular and maintainable than code written using more traditional techniques.\nThe fact that objects are encapsulated produces another important fringe benefit that is sometimes referred to as data abstraction. Abstraction is the process by which complex ideas and structures are made more understandable by the removal of detail and the generalization of their behavior. From a software perspective, abstraction is in many ways the antithesis of hard-coding. Consider a software windowing example: if every detail of every window that appears on a user's screen in a graphical user interface (GUI)-based program had to have all of its state and behavior hard-coded into a program, then both the program and the windows it contains would lose almost all of their flexibility. By abstracting the concept of a window into a window object, object-oriented systems permit the programmer to think only about the specific aspects that make a particular window unique. Behavior shared by all windows, such as the ability to be dragged and moved, can be shared by all window objects.\nThis leads to another basic component of OOT, which is the class. A class includes a set of data attributes plus a set of allowable operations (that is, methods) on the data attributes. Each object is an instance of some class. As a natural outgrowth of encapsulation and abstraction, OOT supports inheritance. A class (called a subclass) may be derived from another class (called a base class, a parent class, etc.) wherein the subclass inherits the data attributes and methods of the base class. The subclass may specialize the base class by adding code which overrides the data and/or methods of the base class, or which adds new data attributes and methods. Thus, inheritance represents a mechanism by which abstractions are made increasingly concrete as subclasses are created for greater levels of specialization. Inheritance is a primary contributor to the increased programmer efficiency provided by OOP. Inheritance makes it possible for developers to minimize the amount of new code they have to write to create applications. By providing a significant portion of the functionality needed for a particular task, classes in the inheritance hierarchy give the programmer a head start to program design and creation. One potential drawback to an object-oriented environment lies in the proliferation of objects that must exhibit behavior which is similar and which one would like to use as a single message name to describe. Consider, for example, an object-oriented graphical environment: if a Draw message is sent to a Rectangle object, the Rectangle object responds by drawing a shape with four sides. A Triangle object, on the other hand, responds by drawing a shape with three sides. Ideally, the object that sends the Draw message remains unaware of either the type of object to which the message is addressed or of how that object that receives the message will draw itself in response. If this ideal can be achieved, then it will be relatively simple to add a new kind of shape later (for example, a hexagon) and leave the code sending the Draw message completely unchanged.\nIn conventional, procedure-oriented languages, such a linguistic approach would wreak havoc. In OOT environments, the concept of polymorphism enables this to be done with impunity. As one consequence, methods can be written that generically tell other objects to do something without requiring the sending object to have any knowledge at all about the way the receiving object will understand the message. Software prograrns, be they object-oriented, procedure-oriented, rule based, etc., almost always interact with the operating system to access the services provided by the operating system. For example, a software program may interact with the operating system in order to access data in memory, to receive information relating to processor faults, to communicate with other processes, or to schedule the execution of a process.\nMost conventional operating systems are procedure-oriented and include native procedural interfaces. Consequently, the services provided by these operating systems can only be accessed by using the procedures defined by their respective procedural interfaces. If a program needs to access a service provided by one of these procedural operating systems, then the program must include a statement to make the appropriate operating system procedure call. This is the case, whether the software program is object-oriented, procedure-oriented, rule based, etc. Thus, conventional operating systems provide procedure-oriented environments in which to develop and execute software. Some of the advantages of OOT are lost when an object-oriented program is developed and executed in a procedure-oriented environment. This is true, since all accesses to the procedural operating system must be implemented using procedure calls defined by the operating system's native procedural interface. Consequently, some of the modularity, maintainability, and reusability advantages associated with object-oriented programs are lost since it is not possible to utilize classes, objects, and other OOT features to their fullest extent possible.\nOne solution to this problem is to develop object-oriented operating systems having native object-oriented interfaces. While this ultimately may be the best solution, it currently is not a practical solution since the resources required to modify all of the major, procedural operating systems would be enormous. Also, such a modification of these procedural operating systems would render useless thousands of procedure-oriented software programs. Therefore, what is needed is a mechanism for enabling an object-oriented application to interact in an object-oriented manner with a procedural operating system having a native procedural interface."}
{"text": "1. Field of the Invention\nThis invention relates to techniques of modulating a carrier by changing its waveform without altering its frequency, phase or amplitude. In particular, this invention relates to a form of modulation called waveform modulation.\n2. Description of the Prior Art\nModulation of a carrier signal with information contained in a modulation signal is well known in the art. Common techniques include pulse modulation, amplitude modulation, phase modulation, frequency modulation and quadrature amplitude modulation. Typically, in such forms of modulation, a sine wave is used as a carrier.\nThe selection of the modulation technique, of course, alters parameters of that carrier sine wave such as whether there is a sine wave (pulse modulation), the amplitude of the sine wave (amplitude modulation), the phase of the sine wave (phase modulation), the frequency of the sine wave (frequency modulation) and the amplitude of both a sine and a cosine wave carrier (quadrature amplitude modulation). Depending upon the selection of the modulation technique, the permitted bandwidth for the trasmitted signal and the expected signal to noise ratio in the communications channel, the maximum data transmission rate can be calculated.\nGenerally, as the sophistication of the modulation technique increases, the complexity of the circuitry required for the transmitter and the receiver increases. For example, in quadrature amplitude modulation, computers are frequently used for demodulation.\nTo maximize the data transmission rate, numerous techniques have been devised for transmitting two signals in the same bandwidth. For example, a common technique used in satellite communications is to have orthogonally polarized waves transmitted at the same frequency. However the polarization of such orthogonally polarized waves can be substantially degraded by precipitation. Other techniques include frequency interleaving such a used in the NTSC television standard.\nTheoretically, for example, a phase shift keyed signal can imposed upon an amplitude modulated carrier to increase the transmission rate. However, the demodulation circuitry for such a transmitted signal is complex to avoid artifacts in the demodulated phase modulated data occurring in the demodulated amplitude signal.\nThus, it is a first goal of this patent to arrive at a new way of modulating a carrier signal with information without requiring complex circuitry for the transmitter or the receiver. It is a further object of the invention to modulate and demodulate the signal in a manner that allows for a maximum data rate. It is yet another object of the instant invention to have the modulated signal to be transparent to other types of signals."}
{"text": "In a general interior permanent magnet motor, as many magnet insertion holes as the number of poles are formed in advance in a circumferential direction at approximately regular intervals through a rotor core formed by laminating and fixing a plurality of magnetic steel sheets punched into a predetermined shape, and one permanent magnet, which has an approximately rectangular cross-sectional shape when viewed in an axial direction, is inserted into each magnet insertion hole. There is a thin portion between both ends of the magnet insertion hole and an outer circumferential surface of the rotor core. This is because a portion between both the ends of the magnet insertion hole and the outer circumferential surface of the rotor core serves as a path for a so-called leakage magnetic flux, which is a magnetic flux flowing out of a surface of the permanent magnet and into another surface of the permanent magnet without passing through a stator core and contributing to an output, and thus this portion is often designed to have a smallest thickness required for securing strength.\nHowever, in the rotor core having such a thin portion, a portion on a radially outer side and a portion on a radially inner side of the magnet insertion hole are connected to each other via the thin portion. Thus, when an upper limit on the number of revolutions of a motor is raised, the strength of the rotor core is required to be improved because the magnitude of a centrifugal force acting on a rotor is proportional to the square of the number of revolutions.\nTo handle such a situation, for example, in Patent Literature 1, there is disclosed a structure in which the magnet insertion hole of one magnetic pole is divided into a plurality of holes in a circumferential direction (thus, a permanent magnet of one magnetic pole is also divided into a plurality of permanent magnets in a circumferential direction), and a bridge configured to connect a portion on a radially outer side and a portion on a radially inner side is formed in that divided position, to thereby intend to achieve improvement of the strength of the rotor core against the centrifugal force. Further, in Patent Literature 1, there is also disclosed forming an arc portion at a corner portion formed at a root of the bridge in consideration of the fact that stress is concentrated on the corner portion formed at the root of the bridge.\nFurther, in Patent Literature 2, there is disclosed forming the bridge itself into a curved shape similarly in consideration of the fact that stress is concentrated on the corner portion formed at the root of the bridge."}
{"text": "Many modern computers rely on various types of components to perform a wide variety of tasks, such as data processing, storage, and communication. Examples of such components include storage drives, network cards, and specialized operation accelerators (e.g., for encryption, media decoding, etc.). If the components of a computer fail to meet the performance requirements of a piece of software, executing that software on the computer may yield an unacceptable user experience. In such a case, the computer user may become frustrated and may mistakenly interpret the software's malfunction as a deficiency of the software, rather than as a deficiency of the computer's other components.\nIn order to avoid or reduce software malfunction and user frustration, software vendors may recommend that their software be executed on computers that are certified to contain components capable of providing specified performance levels. Computer vendors may communicate such certifications to consumers by displaying certification logos on advertisements and packaging for a computer. Accordingly, computer designers who seek to obtain such certifications for their computers may select components based on the performance levels they provide. Accordingly, component designers may design their components to provide the desired levels of performance according to accepted performance metrics. Performance metrics therefore play an important role in both the design and marketing of computers and their components.\nPerformance metrics may be used to standardize measurement and expression of performance levels of computing devices, such as individual components, sets of components, and entire computers. Different performance metrics may relate to different features of a device. For example, some performance metrics may relate to bandwidth, while others may relate to latency. Among the performance metrics that relate to latency, some may describe average latency, while others may describe worst-case latency. Also, the performance of a device may depend on the conditions under which the device's performance is measured. For example, a device that performs well when subjected to one workload may perform poorly when subjected to a different workload.\nMeasuring a device's performance may comprise observing or recording the device's response to a sequence of commands. The sequence of commands may be provided by a synthetic benchmark or a trace. A synthetic benchmark may contain a sequence of commands designed to expose an aspect of the device's performance under stressful conditions, or to emphasize an aspect of the device's performance under idealized conditions. A trace may contain a sequence of commands corresponding to a real-world workload of interest. For example, a trace may contain commands issued to a storage unit by a mobile computer executing a mobile computing workload."}
{"text": "Hospital beds comprise complex mechanical and electronic components for movement, functionality and convenience.\nFoot brakes of prior art hospital beds are typically located on the side under the bed. There are certain disadvantages associated with such foot brakes. For example, during activation, a user such as a nurse has to hold on to the bed, balance on one foot and stretch the other foot under the bed to engage or disengage the brake. As such, if the side rail is in the lower position, visibility is reduced. In addition, if the patient is exiting the bed, the bed may move which is unsafe. Furthermore, the weight of present day beds and patients are relatively large, requiring sufficient braking force to hold a bed in a desired location in a hospital.\nThere is a need for a braking system which is convenient and safe to use. Such a system can be powered in any manner. There is a further need for a braking mechanism that can be manually overridden such as if there is a power failure.\nGenerally, a bed is moved by a series of internal motors and controlled by means of an interface that can be used by users such as hospital personnel or the patient to adjust the bed to suit the comfort and needs of the patient. For safety reasons, the movement of the bed is quite slow and there is a need for an override control, to quickly and efficiently bring the bed into a relatively flat position in case of emergency or for routine tasks such as cleaning, patient transfer or surgery. In past designs, this override function has been initiated through hand controls, foot controls, or a combination of hand and foot controls.\nIn an emergency situation, it is desirable to reposition a bed quickly and easily into a CPR or Trendelenburg position, to facilitate administration of CPR or other resuscitation efforts. The manual- or motor-driven mechanism utilized to raise and lower the Fowler section typically moves too slowly to be acceptable in an emergency situation. Accordingly, emergency releases have been developed to quickly disengage the Fowler section from the drive mechanism to allow for rapid movement, however, these arrangements can be complex, bulky, expensive and difficult to engage and disengage.\nMovement of the foot-end of a hospital bed to various positions that are not aligned with the remainder of the bed, such as a chair position, is difficult when it forms part of the main bed frame\nFor a patient support apparatus in which movement of the Fowler section is effected by a motor-driven mechanism, it would be advantageous to be able to increase the speed at which the Fowler section could be lowered for CPR and Trendelenburg, beyond that speed which is currently obtainable with the motor-driven mechanism powered by a conventional electrical power source.\nEarly designs of adjustable beds often employed the concept of a hand crank and gearing to adjust the height of a bed. Such manual systems suffer from the need for considerable physical effort to adjust the bed height. Other designs include elevation systems incorporating mechanical jacks using hydraulic piston cylinders or screw drives to adjust the height of the hospital bed. Such hydraulic systems are known to be relatively expensive and prone to leakage. Additionally, prior mechanical systems suffer from excessive complexity, excessive size, a lack of load capacity, and manufacturing difficulties.\nHospital bed side rails of the prior art comprise support arms which form undesirable pinch points for users. The movement of such side rails from the deployed to the stowed positions is often hampered by side rail oscillations. The side rail falls due to gravity and the movement can jar the bed and disturb patients.\nIn addition, the patient support apparatus of the prior art relies on batteries to provide all power to the bed's electronic systems. When the battery power runs out, the battery itself must be recharged before power can be supplied to the electronics. This is problematic in circumstances where the life of the battery itself has run out or in settings where a suitable power supply to recharge the battery is not available.\nIn existing apparatuses, the control interface is located on the side or foot-end of a bed. Often, the operator directs movement of the bed from the head-end by pushing on the head-end or push handles located at the head-end. In the event the position of the patient needs to be adjusted while a prior art hospital bed is in motion, the operator has to stop the bed and move around the bed in order to access the bed control interface. If the bed is in a confined space, such as a narrow corridor or elevator, this action may be difficult to execute and result in an undesirable delay in effecting the change in position of the patient.\nCurrently, the angular position of the patient can be determined by measuring the patient's current position with respect to a plane of reference (e.g., the floor or the bed frame). This technique, however, suffers from the drawback that any misalignment in the frame of reference severely affects the integrity of the sensed angular position. Another method for inclinometry is by way of gravitational accelerometers. When the accelerometer is in a stationary position, the only force acting on it is the vertical gravitational force having a constant acceleration. Accordingly, the angular position of the patient can be calculated by measuring the deviation in the inclination angle between the inclination axis and the vertical gravitational force. Although the accelerometers can provide an effective way to measure the inclination in the patient's position, the resolution of the gravitational accelerometers is restricted to a limited range of inclination angles.\nCurrently, nurses and other hospital staff hang pumps (or other hospital equipment) on the top edge of the footboards of hospital beds. Since footboards were not designed to support the hanging of pumps (or other hospital equipment), this current practice reduces access to the controls on footboards, damages foot controls and footboards, generates bed motions and causes damage to pumps (and other equipment) that fall from its hangers.\nOrdinarily, there is a tendency for detached headboards or footboards placed in an upright position against an object or structure to slip, thereby causing the headboard or footboard to fall and potentially suffer damage. This is a particularly acute concern in the situation of a medical emergency during which headboards and footboards may need to be removed and set aside in haste. In a busy hospital, a discarded headboard or footboard that has fallen to the floor creates a tripping hazard to both staff, who may be carrying equipment or medication and thus have an obstructed view of the floor, and patients, who may have compromised mobility owing to illness. Preventing slippage, therefore, reduces the likelihood of personal injury stemming from hastily removed headboards and footboards.\nExisting motorized hospital beds utilize a single speed or multiple defined and preprogrammed speeds for bed movement resulting in the user having to manually switch speeds. Variable speeds in these beds are not automatic\nTherefore, there is a need to provide a patient support apparatus such as a hospital bed which overcomes the problems of the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates in general to the addition of a metal to a superconducting material for homogeneous dispersion of the metal therein, and in particular to a new and improved method for adding silver to a superconducting material to increase the strain tolerance and conductivity thereof in the design and fabrication of magnets and similar products based on the high temperature superconducting material.\n2. Description of the Related Art\nIn the design and fabrication of large scale magnets based on the high temperature superconducting materials, the strain tolerance of the conductor plays an important role. The cuprate superconductors by themselves have a strain tolerance that withstands only about 0.1% to 0.2% strain making fabrication of the magnet exceedingly difficult.\nOne way to fabricate a highly strain tolerant material is to make a microfilamentary composite based on discontinuous filaments. The filaments are separated from one another by a ductile material so that some plastic flow is possible within the matrix without breaking the superconducting filaments. One possible design is to have filaments of superconducting Bi.sub.2 Sr.sub.2 Ca.sub.1 Cu.sub.2 O.sub.8-.delta. in a silver (Ag) matrix. In the design of such a matrix, it is important that the normal barriers between filaments be thin compared to the supercurrent decay length of Cooper paris in silver or about 20 nanometers (nm), see T. Y. Hsiang and D. K. Finnemore, Phys. Rev. B 22,154 (1980) In addition, the length to diameter ratio of the superconducting filaments should be on the order of 10,000 to 1 so that there is enormous overlap area between filaments. As a consequence, the filament-to-filament current density can be low. With this aspect ratio, the conductor will have an effective resistivity in the range of 10.sup.-14 ohms (.OMEGA.) centimeter (cm) even if the silver barriers are totally normal.\nSilver additions to superconducting materials have been demonstrated to have several benefits for improved superconductor performance in addition to promoting strain tolerance in the structure. These benefits appear to be related to the promotion of a controlled crystallographic orientation in the material. This structural improvement can lead to performance enhancements such as high-temperature shift of the zero resistivity temperature (T.sub.C), and an increase in critical current density (J.sub.C), as described in the article by Y. Matsumoto, et al, \"Origin of the Silver Doping Effects on Superconducting Oxide Ceramics\", Appl. Phys. Lett. 56(16), 1585 (1990).\nThe general idea of using silver to provide regions of plastic flow and using overlap areas of the superconducting filaments to overcome the weak link problem seems very promising but it has been difficult to fabricate composites having the desired geometry. Various methods have been used to make Bi.sub.2 Sr.sub.2 Ca.sub.1 Cu.sub.2 O.sub.8-.delta. /Ag composites. In possibly the most straight forward and successful approach, Sato, et al, IEEE Trans. Magn. MAG-27, 1231 (1991), has demonstrated that a multifilamentary conductor composed of 123 monolithic ribbons in a composite degrade only about 30% with a bending strain of about 0.7%. For these samples, there is a high degree of texturing and considerable anisotropy in J.sub.C.\nOther attempts to form high T.sub.C superconducting wire with a high strain tolerance has used the \"powder in a tube\" approach. However, these efforts have met with modest success. In this case, silver was dispersed through the matrix by micro-milling powders followed by composite consolidation using various metal working steps. An important requirement for this technique to be successful is that the particle size of the powders be exceedingly fine. This is necessary so that the subsequent sintering process yields an integral superconducting structure.\nA reliable means of dispersing silver within the superconducting composite has not been described in literature. Current methods are time-consuming involving several process steps to incorporate the silver within the structure. In addition, obtaining a homogenous distribution of silver has proven difficult.\nThus, there is a need for a processing technique for efficient incorporation of silver in a superconducting material to increase the strain tolerance of the conductor to make it suitable for use in large scale magnets."}
{"text": "A gate-array type integrated circuit is a semiconductor device which has basic circuit elements arranged in a pattern on a chip (the \"master slice\"). The circuit elements can be interconnected during manufacture to achieve a desired functionality. Thus, a gate array can be customized.\nCircuit elements typically found in a gate array include storage elements, such as latches, and logic elements, such as AND gates. A \"path\" is defined as a signal transmission passageway between and including circuit elements. There are various types of paths. For example, there are closed paths between storage elements; open paths from a storage element to an output terminal or from an input terminal to a storage element; and through paths from an input terminal to an output terminal. The \"input\" and \"output\" terminals are external terminals (or \"leads\") of the semiconductor circuit device.\nFIG. 11 illustrates sample paths (1) through (4). For example, path (1) includes latches 901 and 902 and gates 909 and 910. Wire portions 917, 918 and 919 connect these elements as illustrated. Each element and wire has an inherent capacitance, which causes a propagation delay as the signal moves from latch 901 to latch 902. These delays are illustrated in the bar graphs of FIGS. 15 and 16. The portion of the propagation delay attributable to each particular element or wire portion is designated in FIGS. 15 and 16 by boxes bearing the reference numeral of the particular element or wire portion.\nAs shown in FIG. 11, latches 901 and 902 receive timing signals T.sub.0 and T.sub.1, respectively. It is essential that the signal propagate from latch 901 to latch 902 during the interval between timing signals (the \"delay reference value\"). Otherwise, the signal will not have reached latch 902 by the time latch 902 is strobbed by timing signal T.sub.1. The delay reference value is shown in FIGS. 15 and 16 by a broken line.\nAs can be seen in FIG. 15, some paths, such as path (4), have propagation delays which exceed the delay reference value. Others, such as path (3), have propagation delays which are less than the delay reference value. As is known, the current between elements can be varied to adjust the propagation delay attributable to the capacitive effects of the wire portions. Specifically, the greater the current, the shorter the propagation delay.\nTo control flow of current between elements, each element includes at its output an emitter-follower circuit, comprised of a plurality of resistors and transistors. As explained in Japanese Post-Examination Application No. 64-4340, these resistors and transistors are selectively coupled to obtain an emitter-follower circuit having the desired current characteristic. By increasing current, the propagation delays can be shortened, so that the aggregated delays for a given path do not exceed the delay reference value.\nThe conventional method of varying an emitter-follower current in accordance with the capacitance of a load circuit will be described with reference to FIGS. 11, 15 and 16. As explained above, FIG. 15 shows the delays in the paths (1)-(4) of FIG. 11 before emitter-follower current is adjusted. The delays illustrated in FIG. 15 include delays in elements and delays on signal lines. FIG. 16 shows the delays after the emitter-follower current has been adjusted in accordance with conventional techniques.\nFor example, consider paths (1) and (2) of FIG. 11. Both paths (1) and (2) have the same number of elements. However, signal line 921 of path (2) is longer than signal lines 917-919 of path (1). Effective capacitance (that is, delay) caused by line 921 is larger than that caused by lines 917-919. Thus, the overall delay of path (2) is greater than the delay of path (1), as shown by the bar graphs of FIG. 15. By increasing the emitter-follower current of a logic element 911, the delay can be shortened, as shown in FIG. 16.\nIn the conventional method, such an adjustment is conducted to all long signal lines, which sometimes causes an unnecessary increase of power consumption. For example, (3) of FIG. 11 includes a long line 923. However, because path (3) has few elements, its overall delay is less than the delay reference time, as shown in FIG. 3. Nevertheless, in accordance with conventional practice, the emitter-follower current of storage element 905 is increased, thus reducing the overall delay of path (3) further below the delay reference time (as shown in FIG. 16), and resulting in additional unnecessary power consumption.\nPath (4) of FIG. 11 includes three elements 914-916 and three short lines 925-928. Because signal lines 925, 926 and 928 are short, their effective capacitances are small. Increasing the emitter-follower current of elements 915 cannot sufficiently reduce the delay of line 927, and it is impracticable to reduce delay on short lines 925, 926 and 928. Therefore, even after adjustment, the delay of path (4) exceeds the delay reference value.\nThus, the conventional technique simply adjusts the emitter-follower current of each element in accordance with the capacitance of the load circuit associated with that element. When the capacitance of the load circuit is large, the element's emitter-follower current is increased, even through the element's path may already have a sufficiently short delay. On the other hand, when the capacitance of the load circuit is small, the current value is not altered, even though the element may be part of a path which has an aggregate delay that exceeds the delay reference value.\nAs described above, the conventional technique reduces the delay between the elements having large capacitive loads. Therefore, the delay of certain individual elements may be reduced. The conventional technique is not effective, however, for reducing the delay in paths which have a larger number of elements connected by short signal lines. Therefore, the speed of the whole semiconductor circuit is limited by the path having the longest delay, and thus limiting the overall speed of the integrated circuit.\nFurther, the conventional method does not reduce the delays attributable to the elements, themselves, as opposed to the delay attributable to the lines between elements. Further, the conventional method does not minimize power consumption."}
{"text": "Nuclear power reactors are well known and are discussed, for example, by M. M. El-Wakil in \"Nuclear Power Engineering\", McGraw-Hill Book Company, Inc., 1962.\nIn a known type of nuclear power reactor, for example, as used in the Dresden Nuclear Power Station near Chicago, Ill., the reactor core is of the heterogenous type. In such reactors the nuclear fuel comprises elongated rods formed of sealed cladding tubes of suitable material, such as a zirconium alloy, containing uranium oxide and/or plutonium oxide as the nuclear fuel, for example, as shown in U.S. Pat. No. 3,365,371. A number of such fuel rods are grouped together and contained in an open-ended tubular flow channel to form a separately removable fuel assembly or bundle as shown, for example, in U.S. Pat. No. 3,431,170. A sufficient number of fuel assemblies are arranged in a matrix, approximating a right circular cylinder, to form the nuclear reactor core capable of self-sustained fission reaction. The core is submerged in a fluid, such as light water, which serves both as a coolant and as a neutron moderator.\nA typical fuel assembly is formed by an array of spaced fuel rods supported between upper and lower tie plates, the rods being several feet in length, on the order of one-half inch in diameter and spaced from one another by a fraction of an inch. To provide proper coolant flow past the fuel rods it is important to maintain the rods in spaced position and restrain them from bowing and vibrating during reactor operation. A plurality of fuel rod spacers spaced along the length of the fuel assembly are provided for this purpose.\nDesign considerations of such fuel rod spacers include the following: retention of rod-to-rod spacing; retention of fuel assembly shape; allowance for fuel rod thermal expansion; restriction of fuel rod vibration; ease of fuel bundle assembly; minimization of contact areas between the spacer and fuel rods; maintenance of structural integrity of the spacer under normal and abnormal (such as seismic) loads; minimization of reactor coolant flow distortion and restriction; maximization of thermal limits; minimization of parasitic neutron absorption; minimization of manufacturing costs including adaptation to automated production. Thus the need to provide such fuel rod spacers creates several significant problems.\nAny material, in addition to the nuclear fuel, that must be used in the construction of the reactor core unproductively absorbs neutrons and thus reduces reactivity with the result that an additional compensating amount of fuel must be provided. The amount of such parasitic neutron absorption is a function of the amount of the non-fuel material, of its neutron absorption characteristics, that is, its neutron absorption cross section, and of the neutron flux density to which it is exposed.\nTo remove the heat from the nuclear fuel, pressurized coolant is forced through the fuel assemblies of the reactor core. The fuel rod spacers in the assemblies act as coolant flow restrictors and cause an undesirable though inevitable coolant flow pressure drop. To maintain proper cooling of the fuel rods along their length and to minimize the required coolant pumping power it is desirable that spacer coolant flow resistance be minimized. The flow resistance of a spacer is a strong function of its projected or \"shadow\" area. Therefore, the flow resistance of a spacer can be minimized by minimizing the projected area of the structure of the spacer. Tests have shown that spacers employing minimized projected area also have the highest thermal limits.\nThe coolant flow resistance of a spacer is also a strong function of the surface or \"wetted\" area of the spacer because of the fluid flow friction between the spacer surfaces and the coolant flowing therethrough. Therefore, the flow resistance of a spacer can be minimized by reducing the height of the spacer.\nAs a practical matter the desire to minimize both parasitic neutron absorption and coolant flow restriction presents a conflict in fuel rod spacer design.\nTo minimize coolant flow restrictions, spacer members should be thin and of minimal cross section area. However, very thin members must be formed of high strength material to provide suitable spacer strength. Also, high strength material with suitable resilience characteristics must be used for any spring member portions of the spacer. It is found that such suitable materials have relatively high neutron absorption characteristics.\nOn the other hand, materials of desirably low neutron absorption characteristics are found to be of relatively low strength, difficult to form and lacking the resiliency desired for the spring member portions of the spacer.\nAn approach toward the resolution of the foregoing design conflict is a \"composite\" spacer wherein the structural members are formed of a material having a low neutron absorption cross section and the spring members thereof are separately formed of suitably resilient material whereby the amount of high neutron absorption cross section material is minimized.\nA variety of such fuel rod spacers have been proposed and used. An example is shown in U.S. Pat. No. 3,654,077. The spacer shown therein (especially the embodiment of FIGS. 5 and 6 thereof) has enjoyed long commercial success. In the spacer thereof the peripheral support member and the divider members are formed of low neutron cross section material such as a zirconium alloy. The divider members are skeletonized, i.e., formed with cutouts, to further reduce neutron loss. To minimize the amount of high neutron cross section spring material in the spacer, a single spring member projects into each of the fuel rod passages, the springs being in the form of four-sided assemblies.\nAnother example of a spacer design is shown in U.S. Pat. No. 3,886,038.\nA further example of a spacer design is a spacer of the ferrule type (a spacer formed of an array of cojoined tubular ferrules) as shown by Matzner et al in U.S. patent application Ser. No. 410,124, filed Aug. 20, 1982, now U.S. Pat. No. 4,508,679 issued April 2, 1985, assigned to the same assignee herein, which application is hereby incorporated by reference herein.\nIn the past, fuel assemblies were designed for a residence time in the core of in the order of four years. Recent trends toward longer fuel burnup require fuel assembly residence times in the order of six years or more.\nThis increased residence time in the core gives rise to a further spacer design problem, namely, an increase in the amount of hydrogen picked up by the spacer from the environment of the core to the extent that the hydride concentration in the material of the spacer can cause embrittlement thereof and consequent decrease in its strength. If the hydride concentration becomes too high, there is a possibility of spacer failure.\nThe amount of hydrogen picked up by the spacer is proportional to the spacer surface area which is exposed to the coolant, i.e., the \"wetted\" surface area. The rate of hydrogen diffusion into the relatively thin spacer members is sufficient to give a substantially uniform hydride concentration throughout the volume of the spacer material. Therefore, the hydride concentration is proportional to the ratio of the wetted surface area to the volume of the spacer material.\nThus the hydride concentration can be reduced by increasing the cross section area of the spacer members by an increase in their width or thickness.\nAs discussed hereinbefore, coolant flow resistance through a spacer is a function of both cross section area and height of the spacer members, therefore, to maintain a desirably low coolant flow resistance, any increase in thickness of the spacer members must be compensated for by a decrease in the height of the spacer members.\nThus, for minimizing hydride concentration, a spacer of minimum height is indicated.\nIn many spacer designs the minimum height is limited by the design of the spacer springs. In spacers where the springs are vertically (axially) oriented, such as shown in the above-mentioned patent application Ser. No. 410,124, minimum height of the spacer is dictated by the spring length needed for the required spring flexibility and force.\nAn object of the invention is a nuclear fuel element spacer with spacer members having sufficient cross section area to maintain hydride concentrations therein at suitable levels over extended residence time of the spacer in a nuclear reactor core.\nAnother object is to minimize the height of a spacer to minimize the resistance to coolant flow therethrough.\nAnother object is a laterally or horizontally oriented spring member for a spacer so that the required spring length does not dictate the minimum height of the spacer.\nAnother object is a lateral spacer spring which spans two adjacent fuel element passages for lateral support of the elements extending through the passages.\nAnother object is a lateral spacer spring which is retained by horizontal slots in the spacer members whereby a substantial portion of the spring is within the \"shadow\" of the spacer members so that the contribution of the spring to coolant flow resistance is minimized.\nAnother object is a lateral spacer spring wherein the stress distribution is such as to make efficient use of the spring material whereby the amount of spring material in the spacer is minimized.\nAnother object is a spacer formed of an array of cojoined ferrules wherein the ferrules have outer surfaces of octagonal shape and inner surfaces of circular shape."}
{"text": "The use of an assembly comprising a ceramic head and a metal femoral rod as a femoral prosthesis is known, one end of the rod being implanted in the femur and the other being nested in a blind recess cut into in the ceramic head. This assembly in general is associated with a natural or artificial acetabular cup supporting the rotation of the ceramic head in order to form an internal prosthesis for the femoral joint.\nThe blind recess made in the ceramic head (the female part) is customarily of truncated cone shape with a circular cross section and slight conicity; the male part of the femoral metal rod which is to be nested in it has practically the same angle at the apex, so that assembly by nesting the parts together is done without play between the two pieces, the relative rotation of the head on the rod being prevented by various means such as jointing by force, use of adhesives, wedging, brazing, etc.\nBecause the angle at the apex of the cone is slight, it quite often occurs that during the placement of the two pieces by the physician their alignment is not totally correct; thus instead of having a contact surface between the faces of the female and male pieces, the male piece is actually only in contact with some points on the surface of the recess of the female piece. For the same reasons it also happens that the male piece is not nested in place over all of the bearing surface provided in the recess of the female piece; for example, the apex of the male piece might be incorrectly supported by a point on the surface of the recess which is not located in the bottom of the recess and, at the same time, the contact surface of the male piece might be incorrectly supported on a point of the recess close to the periphery of the ceramic head.\nSometimes the male metal cone and the ceramic recess also have faults in their geometry which cause the contact between the two pieces to occur only at some points so that their placement is not correct.\nAll of these problems, associated with the placing of the male and female pieces by means of nesting them together, by definition mean that the mechanical stresses due to exterior stress and sustained by the ceramic head are transmitted from the male piece to the femur by means of only a few contact points between the male piece and the recess of the female piece, instead of by means of a contact surface. Thus the stresses supported by the ceramic head are considerably increased at these points and are the cause of a very serious increase of a risk of breaking of the ceramic head. These breaks usually occur in the anterior area, rather than the deep area, of the recess of the ceramic head."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of nucleic acid sequence detection and a nucleic acid sequence detection substrate. Particularly, the present invention relates to a technique for highly sensitively detecting, conveniently and at low cost, the presence of a particular target nucleic acid sequence to be detected in a sample.\n2. Description of the Related Art\nA conventional method for detecting a nucleic acid sequence specific to a sample as a target involves: binding the target nucleic acid sequence to oligonucleotide strands immobilized on a substrate and binding a different oligonucleotide strand labeled with a fluorescent substance to the bound target nucleic acid sequence; and detecting the presence of the particular nucleic acid sequence in the sample by detecting/amplifying a reflected light signal from the fluorescent substance label by irradiation with light (e.g., laser beam) at a particular wavelength. Another detection method involves: binding the target nucleic acid sequence to oligonucleotide strands immobilized on a substrate; performing elongation reaction using the oligonucleotide strands as primers and the bound target nucleic acid sequence as a template, wherein base substances incorporated during the reaction are labeled with a fluorescent substance; and detecting/amplifying a reflected light signal from the fluorescent substance label by irradiation with light (e.g., laser beam) at a particular wavelength.\nHowever, these detection methods must detect a fluorescent substance and therefore require a special detection apparatus that emits light at a particular wavelength. Such a detection apparatus is expensive and has a limited use to some research institutes or universities under the present circumstances. Thus, these approaches cannot be adopted, for example, for conveniently detecting pathogenic microorganisms such as norovirus or Cryptosporidium in water (sample) at water treatment sites, due to complicated procedures and expensive analysis apparatuses.\nAgainst this backdrop, for example, the MPEX (Multiple Primer EXtension) method described in K. Kinoshita et al., Multiple primer extension by DNA polymerase on a novel plastic DNA array coated with a biocompatible polymer, Nucleic Acid Research, Vol. 35, No. 1, 2007, pp. e3, and Japanese Patent Application Laid-Open Nos. 2006-174788 and 2007-222010 has been developed as a method for conveniently detecting a particular nucleic acid sequence at detection (e.g., water treatment sites) at low cost.\nThis MPEX method is a genetic testing method using DNA elongation reaction catalyzed by an enzyme (DNA polymerase) on a substrate and is also used in gene mutation analysis, SNP (Single Nucleotide Polymorphism) analysis, and microorganism identification. This method uses, as primers, oligonucleotide strands immobilized on a substrate having, on the surface thereof, a polymer substance containing a first unit having a group induced from phosphoric ester constituting the hydrophilic moiety of phospholipid and a second unit having an active ester group. It involves: after addition of a target nucleic acid sequence, heat-denaturing the double-stranded DNA into single strands on the substrate and performing DNA elongation reaction catalyzed by DNA polymerase using the single strand as a template to incorporate bases, all or some of which are modified with a ligand, into the elongation reaction product (amplified DNA); and finally adding thereto an active substance reactive to the ligand for reaction. The MPEX method thus achieves highly sensitive gene detection."}
{"text": "The present invention relates to disk drive data storage devices, and in particular to disk and hub assemblies used in disk drive storage devices.\nThe extensive data storage needs of modern computer systems require large capacity mass data storage devices. A common storage device is the rotating magnetic disk drive.\nA disk drive typically contains one or more smooth, flat disks which are rigidly attached to a common spindle. The disks are stacked on the spindle parallel to each other and spaced apart so that they do not touch. The disks and spindle are rotated in unison at a constant speed by a spindle motor.\nEach disk is formed of a solid disk-shaped base or substrate, having a hole in the middle for the spindle. The substrate is commonly aluminum, although glass, ceramic, plastic or other materials are possible. The substrate is coated with a thin layer of magnetizable material, and may additionally be coated with a protective layer.\nData is recorded on the surfaces of the disks in the magnetizable layer. To do this, minute magnetized patterns representing the data are formed in the magnetizable layer. The data patterns are usually arranged in circular concentric tracks. Each track is further divided into a number of sectors. Each sector thus forms an arc, all the sectors of a track completing a circle.\nA moveable actuator positions a transducer head adjacent the data on the surface to read or write data. The actuator may be likened to the tone arm of a phonograph player, and the head to the playing needle.\nThere is one transducer head for each disk surface containing data. The transducer head is an aerodynamically shaped block of material (usually ceramic) on which is mounted a magnetic read/write transducer. The block, or slider, flies above the surface of the disk at an extremely small distance as the disk rotates. The close proximity to the disk surface is critical in enabling the transducer to read from or write to the data patterns in the magnetizable layer. Several different transducer designs are used, and in some cases the read transducer is separate from the write transducer.\nThe actuator usually pivots about an axis to position the head. It typically includes a solid block near the axis having comb-like arms extending toward the disk, a set of thin suspensions attached to the arms, and an electromagnetic motor on the opposite side of the axis. The transducer heads are attached to the suspensions, one head for each suspension. The actuator motor rotates the actuator to position the head over a desired data track. Once the head is positioned over the track, the constant rotation of the disk will eventually bring the desired sector adjacent the head, and the data can then be read or written.\nAs computer systems have become more powerful, faster, and more reliable, there has been a corresponding increase in demand for improved storage devices. These desired improvements take several forms. It is desirable to reduce cost, to increase data capacity, to increase the speed at which the drives operate, to reduce the electrical power consumed by the drives, and to increase the resilience of the drives in the presence of mechanical shock and other disturbances.\nIn particular, there is a demand to reduce the physical size of disk drives. To some degree, reduction in size may serve to further some of the above goals. But at the same time, reduced size of disk drives is desirable in and of itself. Reduced size makes it practical to include magnetic disk drives in a range of portable applications, such as laptop computers, mobile pagers, and xe2x80x9csmart cardsxe2x80x9d.\nAn example of size reduction is the application of the PCMCIA Type II standard to disk drives. This standard was originally intended for semiconductor plug-in devices. With improvements to miniaturization technology, it will be possible to construct disk drives conforming to the PCMCIA Type It standard.\nIn order to shrink the size of disk drives, every component must be reduced in size as much as possible. Additionally, new designs of existing components must be developed to permit reduced size and make assembly of miniaturized components practical.\nOne limit on the extent to which a disk drive can be reduced in size is the disk/hub assembly. A conventional disk/hub assembly comprises a cylindrical hub having a flange at the bottom for supporting a disk stack. A motor for rotating the disks is located within the hub. The stack of disks rests on the upper surface of the flange, the hub fitting within corresponding holes of the disks. The individual disks are separated by spacer rings surrounding and adjacent to the hub. A clamping apparatus is attached to the top of the hub and applies a downward force to the disk stack, forcing the bottom disk against the flange and holding the stack in place. Where a disk drive has only a single disk, essentially the same design is used, but the clamping apparatus clamps only the one disk instead of the disk stack. The clamping apparatus is typically a flat steel ring having a formed circular ridge near its outer edge. The flat portion of the ring is attached to the upper surface of the hub with screws, while the ridge portion applies pressure to the disk stack or single disk. Several alternative clamp designs exist, but all involve multiple parts.\nThe conventional disk drive disk/hub design is not well suited to very small form factor disks, such as the PCMCIA Type II form factor. The hub must be sufficiently large to accommodate the screws. Even where very small screws are used, this requirement adds to the size and weight of the hub. The extremely small parts make assembly difficult. There is difficultly tolerating high mechanical shock of portable applications. The relatively thin disks tend to warp when clamped with sufficient clamping force; even a small warpage can be serious when track widths are being reduced. Finally, although the drive is much smaller, simple scaling down of the size of conventional parts will not result in any significant cost reduction; in fact, it may increase costs. It is desirable to develop an alternative hub/disk assembly which reduces costs and is more suitable to the design requirements of small form factor disk drives.\nIt is, therefore, an object of the present invention to provide an enhanced disk drive storage apparatus.\nAnother object of this invention is to reduce the cost of a disk drive storage apparatus.\nAnother object of this invention is to provide an enhanced hub/disk assembly for a small form factor disk drive storage device.\nAnother object of this invention is to provide a reduced cost hub/disk assembly for a small form factor disk drive.\nAnother object of this invention is to provide a hub/disk assembly having fewer parts.\nAnother object of this invention is to provide a hub/disk assembly which is easier to fabricate and assemble.\nAnother object of this invention is to provide a hub/disk assembly having greater resistance to mechanical shock.\nAnother object of this invention is to provide a hub/disk assembly having reduced distortion of the disk.\nA disk for a disk drive is mounted on an integrally-formed one-piece hub, clamp and rotor housing. The hub is preferably formed of molded plastic. In the preferred embodiment, the hub includes a substantially hollow cylindrical portion for surrounding a set of bearings mounted on a central shaft, and a flange portion extending from the cylindrical portion near the middle of the cylinder. The flange supports a single disk, which rests on its upper surface. A plurality of circumferentially spaced compliant clamping fingers extend upward from the upper surface of the flange to engage the disk through its central hole. The fingers are inclined slightly outward, and include a pawl at the end of each finger to engage the disk. The edges of the central hole of the disk are preferably chamfered to improve engagement with the fingers and center the disk. To assemble the disk to the hub, the disk is simply forced over the pawls of the fingers until they engage the chamfered edge. The fingers, being slightly elastic, will deform inward sufficiently to permit the disk to slide over the pawls, and will then return outward to lock the disk in place. The fingers automatically center the disk about the hub axis, without the need for centering tools.\nThe hub also serves as a housing for the rotor of a spindle motor which rotates the disk. In the preferred embodiment, a brushless DC spindle motor is located under the flange. A back iron ring and permanent magnets are attached to the underside of the flange at its outer edge to form the rotor portion of the motor. An electromagnetic stator assembly surrounds the disk axis in the annular space defined by the flange, the permanent magnets and back iron at the outer edge, and the cylindrical portion of the hub at the inner edge."}
{"text": "The invention pertains to methods of forming polished material, such as, for example, methods of forming isolation regions.\nIn modern semiconductor device applications, millions of individual devices are packed onto a single small area of a semiconductor substrate, and many of these individual devices may need to be electrically isolated from one another. One method of accomplishing such isolation is to form a trenched isolation region between adjacent devices. Such trenched isolation region will generally comprise a trench or cavity formed within the substrate and filled with an insulative material, such as silicon dioxide.\nA prior art method for forming trench isolation regions is described with reference to FIGS. 1-10. Referring to FIG. 1, a semiconductive wafer fragment 10 is shown at a preliminary stage of a prior art processing sequence. Wafer fragment 10 comprises a semiconductive material 12 having an upper surface 13. A layer of silicon dioxide 14 is formed over upper surface 13, and a layer of silicon nitride 16 is formed over silicon dioxide 14. A patterned masking layer 18 is formed over silicon nitride 16.\nSubstrate 12 can comprise, for example, monocrystalline silicon lightly doped with a conductivity enhancing material. For purposes of interpreting this document and the claims that follow, the term xe2x80x9csemiconductive substratexe2x80x9d is defined to mean any construction comprising semiconductive material, including, but not limited to, bulk semiconductive materials such as a semiconductive wafer (either alone or in assemblies comprising other materials thereon), and semiconductive material layers (either alone or in assemblies comprising other materials). The term xe2x80x9csubstratexe2x80x9d refers to any supporting structure, including, but not limited to, the semiconductive substrates described above.\nSilicon dioxide layer 14 has a typical thickness of about 90 xc3x85, nitride layer 16 has a typical thickness of from about 700 xc3x85 to about 800 xc3x85, and masking layer 18 has a typical thickness of about 10,000 xc3x85. Nitride layer 16 comprises a lower surface 15 and an upper surface 17. Patterned masking layer 18 can comprise, for example, photoresist.\nReferring to FIG. 2, a pattern is transferred from layer 18 to nitride layer 16 and oxide layer 14 to form masking blocks 20 over substrate 12. Blocks 20 are separated by intervening openings (also referred to as trenches or gaps) 22, 24, 26, 28 and 30. Gaps 22, 24, 26, 28 and 30 define locations wherein isolation regions will ultimately be formed. In the shown typical embodiment, gaps 22, 24, 26, 28 and 30 vary in width, with gap 26 representing a wide trench. The variation in width of gaps 22, 24, 26, 28 and 30 is a matter of discretion for persons fabricating devices over substrate 12. Such variation in width is shown to exemplify how particular problems associated with polishing processes (the problems are discussed below) become exaggerated at wider trench openings.\nReferring to FIG. 3, openings 22, 24, 26, 28 and 30 are extended into substrate 12. The processing of FIGS. 2 and 3 typically occurs in a single etch step.\nReferring to FIG. 4, photoresist material 18 (FIG. 3) is removed. Subsequently, an insulative material 40 is provided to fill the trenches 22, 24, 26, 28 and 30. Insulative material 40 typically comprises silicon dioxide, and can be formed by, for example, high density plasma deposition. The term xe2x80x9chigh densityxe2x80x9d as used in this document to refer to a deposition plasma means a deposition plasma having a density of greater than 1010 ions/cm3. The outer surface of the deposited insulative material 40 comprises surface peaks 42 (which actually have a three-dimensional pyramid shape) extending over the patterned nitride material 16. Peaks 42 result as an aspect of high density plasma deposition.\nReferring to FIG. 5, wafer fragment 10 is subjected to a polishing process to remove insulative material 40 from over nitride layer 16, and to thereby remove peaks 42 (FIG. 4) and polish the insulative material 40 down to about even with an upper surface of nitride 16. After the polishing process, insulative material 40 forms isolation regions 50, 52, 54, 56 and 58 within openings 22, 24, 26, 28 and 30, respectively. The isolation regions and silicon nitride layer 16 comprise a coextensive upper surface 60. Upper surface 60 can be at about a same elevational level as original upper surface 17 (FIG. 1) of nitride layer 16, or can be below such elevational level, depending on whether the polishing process has removed any of the material of silicon nitride layer 16. Generally, the polishing process removes some of silicon nitride layer 16, but will slow significantly upon reaching silicon nitride layer 16 such that silicon nitride layer 16 effectively functions as an etch stop layer to define an end point of the polishing process. Exemplary polishing processes include chemical-mechanical polishing, as well as dry and wet etches selective for silicon dioxide relative to silicon nitride. In particular applications, insulative material 40 comprises silicon dioxide and is polished down to a level that is above the upper surface of the silicon nitride. The material 40 is then brought to about level with the upper surface of the silicon nitride with a subsequent wet acid (hydrofluoric acid) dip.\nIdeally, upper surface 60 will be substantially planar. However, as shown a problem frequently occurs during the polishing of insulative material 40 wherein concavities 62 occur at the top of isolation regions 50, 52, 54, 56 and 58. Such problem is commonly referred to as xe2x80x9cdishing.xe2x80x9d The problem is frequently more severe at wider isolation regions, corresponding to wider trenches, (i.e., isolation region 54) than at narrower isolation regions (i.e., isolation regions 50, 52, 56 and 58). The dishing can become particularly pronounced for isolation regions having widths greater than or equal to about 5 microns.\nReferring to FIGS. 6 and 7, nitride layer 16 (FIG. 5) is removed, and subsequently pad oxide 14 is stripped. The stripping of pad oxide layer 14 can be accomplished with a hydrofluoric acid dip. Such dip also removes some of the silicon oxide from isolation regions 50, 52, 54, 56, and 58, and accordingly reduces a height of the oxide in such isolation regions.\nReferring to FIG. 8, a sacrificial silicon dioxide layer 19 is formed over substrate 12. Such sacrificial silicon dioxide layer can be grown from silicon of substrate 12, or deposited over substrate 12.\nReferring to FIG. 9, sacrificial silicon dioxide layer 19 is stripped from over substrate 12. The stripping of sacrificial silicon dioxide layer 19 can be accomplished with a hydrofluoric acid dip. Such dip also removes some of the silicon oxide from isolation regions 50, 52, 54, 56, and 58, and accordingly reduces a height of the oxide in such isolation regions.\nReferring to FIG. 10, a gate oxide layer 21 is formed over substrate 12. Such gate oxide layer can be grown from silicon of substrate 12, or deposited over substrate 12. It is noted that only pertinent processing steps are discussed in describing FIGS. 8-10, and that additional processing steps (in addition to those discussed) can occur between the forming of the sacrificial oxide layer and the forming of the gate oxide layer.\nAs shown in FIG. 10, the dishing described with reference to FIG. 5 can remain in isolation regions 50, 52, 54, 56 and 58 after the processing of FIGS. 6-10, and can result in at least some of the isolation regions having portions below an elevational level of the upper surface of substrate 12. The dishing also causes the isolation regions to have corners 70 which are not right angles. Such corners 70 can undesirably affect operating voltages of devices formed proximate isolation regions 50, 52, 54, 56 and 58, and can, for example, result in non-uniformity of threshold voltages for such devices. It would therefore be desirable to develop alternative methods of forming isolation regions.\nIn one aspect, the invention encompasses a method of forming a polished material. A substrate is provided and an elevational step is provided relative to the substrate. The elevational step has an uppermost surface. A material is formed beside the elevational step. The material extends to above the elevational step uppermost surface and has lower and upper layers. The lower layer polishes at slower rate than the upper layer under common polishing conditions. The lower layer joins the upper layer at an interface. The material is polished down to about the elevational level of the elevational step uppermost surface utilizing the common polishing conditions.\nIn another aspect, the invention encompasses a method of forming an isolation region. A substrate is provided. The substrate has an opening extending therein and a surface proximate the opening. A material is formed within the opening. The material extends to above the substrate surface, and comprises a lower layer and an upper layer. The lower layer is more dense than the upper layer, and joins the upper layer at an interface that extends to at or below an elevational level of the substrate surface. The material is polished at least down to about the elevational level of the substrate surface."}
{"text": "The Internet of Things (IoT, also Cloud of Things or CoT) refers to the interconnection of uniquely identifiable embedded computing-like devices within the existing Internet infrastructure. Typically, IoT is expected to offer advanced connectivity of devices, systems, and services that goes beyond machine-to-machine communications (M2M) and covers a variety of protocols, domains, and applications. The interconnection of these embedded devices (including smart objects), is expected to usher in automation in nearly all fields, while also enabling advanced applications like a Smart Grid. These connected devices could be within a home (e.g. thermostat or garage door opener), a workplace (e.g. building thermostat or lights), or within a factory (gas sensors or machines).\nAccording to Gartner, there will be nearly 26 billion devices on the IoT by 2020. ABI Research estimates that more than 30 billion devices will be wirelessly connected to the IoT by 2020. There are cases in which the information (data packets) sourced from electronic internet connected devices are consumed within a zone and/or a geographical area. The zone could be inside a building, a campus and/or a geographical area. In most cases, information provided by electronic internet connected devices needs to be consumed within a zone of interest. The security of the devices becomes a major issue when information provided with these devices could be harvested for malicious activities outside of the zone of interest. For example, a thermostat or metering devices connected to the internet may provide information to outsiders about when the households are home or away.\nIn order to solve this problem, it is advantageous to provide a secure network architecture to accommodate internet connected devices."}
{"text": "The ability to conduct high-speed data and voice communications between remotely separated data processing systems and associated subsystems has become a requirement of a variety of industries and applications, such as business, educational, medical, financial and personal computer uses. Moreover, it can be expected that future applications of such communications will engender more systems and services in this technology. Associated with such applications has been the growing use and popularity of the “Internet”, which continues to stimulate research and development of advanced data communications systems between remotely located computers, especially communications capable of achieving relatively high-speed data rates over an existing signal transport infrastructure (e.g., legacy copper cable plant).\nOne technology that has gained particular interest in the telecommunication community is digital subscriber line (DSL) service, which enables a public service telephone network (PSTN) to deliver relatively high bandwidth signals (including voice and data) using conventional telephone company copper wiring infrastructure. DSL service has been categorized into several different technologies, based upon expected data transmission rate, the type and length of data transport medium, and encoding/decoding schemes.\nRegardless of its application, the general architecture of a DSL network essentially corresponds to that diagrammatically shown in FIG. 1 as comprising a pair of remotely separated, mutually compatible digital communication transceiver entities. One entity is located at respective network controller site 10 (such as a telephone company central office (CO)), while a second entity is located at a customer premises site 20. Each transceiver is coupled to a communication link, such as a twisted pair (loop) 30 of an existing copper plant. Using ATM-based, digital subscriber line (DSL) protocol, this telecommunication fabric allows information, such as voice and (Internet-sourced) data (which is readily accessible via a backbone network 15), to be transmitted from the central office site 10 over the DSL loop 30 to an integrated access device (LAD) 21 at the customer site 20.\nFor this purpose, at the network controller site 10, a DSL transceiver 11 is customarily located in a DSL access multiplexer (DSLAM) 12. Within the communication infrastructure of the telephone company, DSLAM 12 is coupled with the backbone 15, which typically contains one or more of signaling transport devices, such as an asynchronous transfer mode (ATM) switch 31, a voice gateway 33, a Class-5 switch 35, and the like, that are linked to an internet service provider (ISP) 37. Also a data gateway 36 may link the ATM switch 31 to a data network 38.\nThe other transceiver, serving the customer premises site 20, may comprise an integrated access device (IAD) 21, which is coupled via a plain old telephone system (POTS) interface 23 to a modem 25 (such as a V.90 modem) serving data terminal equipment (DTE) 27.\nFor transporting data and voice, an ATM network of the type shown in FIG. 1 employs ATM Adaptation Layer 5 (AAL5) for data transport, and AAL2 for voice transport. As ATM is a ‘cell’-based asynchronous transfer protocol, processing at both is the transmit or source site and the receiver entity are necessary to ensure a continuous CBR flow of voice and voice-band data cells across the ATM fabric. Unfortunately, AAL2 protocol-based voice and voice-band data transmission can be disrupted by delays encountered by the ATM cells during their transport over the network. These delays are of two types: 1—fixed delay associated with the configuration of the network (which is predictable and readily accounted for), and 2—variable delay (termed Cell Delay Variations (CDV)) associated with the traffic load on the network switches, causing successive ATM cells to arrive at a receiving or destination entity in an aperiodic manner.\nTo minimize or eliminate these disruptions, in order to effectively ensure reliable voice-band data transmission, it is necessary to remove the variable delay component of cell arrival time. This is customarily achieved through the use of a cell jitter buffer of sufficient length to accommodate maximum cell delay variation), and synchronizing the receive site's POTS interface (the IAD's CODEC) clock to the far-end or source site's transmitter clock, in a manner that avoids overflow or underflow of the buffer (which will occur if the clocks are not locked together).\nOne relatively straightforward method to recover the clock is to encode the transported ATM stream with a Synchronous Residual Time Stamp (SRTS) representative of the frequency difference between the source clock and a common reference network clock. At the receiving entity, the SRTS is decoded to regenerate the source clock frequency. Unfortunately, for AAL2-based data transmissions, physical layer timing on the DSL loop may not always be traceable to a primary reference source (transmit site) clock, and must be extracted ‘adaptively’ from the incoming AAL2 cell stream.\nIn an adaptive clock recovery scheme, no explicit timing information is transmitted from source to destination across the network and no common reference clock is used. Instead, source clock frequency information is derived by monitoring ATM cell arrival activity, and averaging out CDV effects. While there is currently no ‘standardized’ method, adaptive clock recovery has typically involved monitoring the ‘fill’ level of a cell jitter buffer, through which received ATM cells are controllably clocked by an associated clock recovery loop, and adjusting the receive entity clock, so that positions of write/read pointers to the buffer fall within a prescribed error window relative to a selected (e.g. median or statistically averaged) buffer fill level, and avoid overflow or underflow of the buffer.\nFor an illustration of non-limiting examples of literature describing various clock recovery schemes including both SRTS and buffer fill level-based adaptive mechanisms of the type described above, attention may be directed to the following U.S. Pat. Nos. 5,361,261, 5,844,891, 5,966,387, 6,111,878, 6,188,692 and 6,252,850."}
{"text": "This invention relates to a temperature responsive electric fuse or switch of the non-reset type, operable to open a circuit at a specific temperature. More particularly, the invention relates to an improvement of a temperature sensitive pellet for a thermal fuse, the reliability of which is improved by a sensitive and rapid temperature response characteristic at a specific temperature.\nIt is known to use overheating prevention devices to open a circuit when the temperature of an electric apparatus exceeds a given range, to increase the safety of the apparatus. Conventionally, such overheating preventing devices include two well-known types. One type is the non-reset type thermal fuse, employing temperature sensitive materials fusible at a specific temperature. Another type is self-resetting and uses bimetal thermocouple means. The bimetal-type is disadvantageous, in that even when it once functions to open the circuit, a drop in the ambient temperature allows the device to reset and start the current flow again. The switch-off and switch-on operation may thus be repeated unless the power supply is cut off, or the cause of the fault is eliminated. Thus, heat may be gradually stored in the bimetal device, whereby the temperature response of the device may be changed.\nRecently, the non-reset type thermal fuse which employs an organic temperature sensitive pellet has been widely used due to the advantage of a constantly stable function over long periods of time in the absence of external changes in comparison with thermal fuses employing fusible metals. Such organic pellet fuses do not store any heat due to its non-reset characteristic, whereby the desired safety results.\nThis kind of thermal fuse of the non-reset type is described, for instance, in U.S. Pat. No. 3,519,972. FIG. 1 shows such a conventional thermal fuse with a cylindrical housing 10, made of a metal having a good electrical and heat conductivity. A first lead-in wire 12 is connected to one end of the housing 10. The switching elements include a temperature sensitive pellet 14, metallic disks or plates 16 and 18, two types of compression spring means 20 and 34 arranged within the housing 10, and a resilient slidable contact member 22, the peripheral portion 24 of which abuts against the inner wall of the metallic housing 10. These elements and an insulator 28 are enclosed by means of a housing edge 26. A second lead-in wire 30 passes through a further insulator 36 which hermetically seals the housing. Two different compression springs 20 and 34 are combined so that one of the springs has a greater resiliency than the other, while the other spring has a greater restoring force than the former. The second lead-in wire 30 has a head portion 32 forming a contact. A more detailed description of the prior art fuse and its function may be had from the above patent.\nIf the temperature at which the fuse is supposed to operate is, for example, 130.degree. C, such a substance as fructose is used for the pellet. If the temperature is 145.degree. C, glucose is used and if the temperature is 155.degree. C, salicylic acid is used for the pellet, to give but a few examples. The organic substance, which is used to prepare the temperature sensitive pellets, is first pulverized into a powder which is then molded into pellets by a press molding.\nHowever, the thermal conductivity of the pellel 14 is not always good. Moreover, the transmission of thermal energy always occurs from the outside of the pellet inwardly and through the cylindrical housing 10. Thus, even if the outside of the pellet should melt, its center may remain solid. As a result, it happens sometimes that the slidable contact member is not separated from the head portion of the second lead-in wire and the electric circuit is not interrupted. Moreover, since the melting of the temperature sensitive pellet occurs gradually from the outside inwardly, the slidable contact member is in an unstable state during the meeting. Thus, even if the slidable contact member is still separated from the second lead-in wire, the distance between the contacts is small enough to cause sparks, and the current is never fully interrupted. This is a serious disadvantage of the conventional temprature temperature pellet.\nThe conventional temperature sensitive pellet of the above fuse is made of an insulating organic material and molded under the proper pressure to form, for example, a cylinder of fine grains of said organic materials which are disclosed in U.S. Pat. No. 2,934,628. This kind of material must be selected in accordance with the particular use for which the thermal fuse is intended. Besides, it is difficult to reduce the costs of this prior art, small-sized thermal fuse, because it is hard to prepare the selected fine grains of organic materials."}
{"text": "The present invention relates to an adjustable supporting device, especially intended for premature infants, that is, infants who are born too soon, in open and closed incubators.\nHowever, the subject matter of the invention may also be used both in cots in premature wards in hospital and also in other wards such as medical, surgical and neurosurgical wards where there are infant patients, and the subject matter of the invention will also be suitable for home use.\nThe invention is based on the principle of separate inflatable sections or units wherein in addition the pressure can be adjusted individually.\nPremature infants in need of incubator treatment belong to a highly vulnerable group of patients in need of extremely careful treatment.\nA feature of this careful treatment is that the infants should be xe2x80x9chandledxe2x80x9d as little as possible when they are in an incubator or a cot.\nThe object of the subject matter of the present invention, an adjustable supporting device for such infants, is to solve this problem by providing a supporting device divided into at least two sections which can be inflated and pressure-adjusted individually, thereby allowing the infant\"\"s position to be changed without any direct manual contact.\nA great number of devices are known, all of which aim to have a certain form of individual inflatability and pressure adjustment.\nIn this connection reference will made to DE-AS 1 012 737 which describes an air mattress of a known type per se, but wherein the tubular bodies which constitute the side edges and foot piece can be raised just above the vertical plane and joined together while obtaining a bed-like structure.\nDE-PS 918 109 describes a mattress wherein the middle, inflatable sections can be turned 90xc2x0 upwards relative to the outer inflatable section and wherein optionally there may in addition be placed inflatable cushions between the thus produced walls.\nNO 129 279 describes a mattress-like structure which with the aid of transverse and longitudinal seams is divided into sections, and wherein side sections and end sections can be folded up to form a bed when the air is released from the sections located in the corners.\nGB 1 602 952 described a supporting device consisting of a number of parallel, inflatable bodies, intended to extend transverse of a person resting on the supporting device and wherein each of these bodies can be filled with a fluid, and wherein each body has an upper part which can be emptied separately to provide an certain adjustment to the body contours.\nGB 1 534 821 describes a device that is in principle the same as that in the aforementioned \"\"952, but used in connection with a dentist\"\"s chair or the like.\nFinally, NO 174 452 describes an individually adjustable cushion system wherein separately manipulatable cushions can via a manifold be brought to the desired pressure for placing individually under different parts of a patient.\nThe object of the subject matter of the present invention is to simplify the prior art whilst aiming to satisfy the needs of premature infants in connection with a supporting device.\nAs mentioned, the subject matter of the invention is intended to be used in incubators, but can of course be used in any other place, including domiciliary care.\nAccordingly, the present invention relates to an adjustable supporting device, chiefly for premature infants, comprising an inflatable base, consisting of two superimposed air and gas impermeable layers, divided into at least two flat sections, and a circumferential, inflatable wall, optionally also divided into at least two sections, wherein all the base sections and wall or wall sections are inflatable and pressure-adjustable independent of one another.\nPreferably, the wall is in one section and is in the shape of a circular sausage that forms a complete circle.\nThe premature supporting device of the invention preferably has a basic shape that is generally oval.\nAlthough the base may be divided into any number of flat sections, it is in principle preferred that there be only two sections.\nIn order to prevent the infant lying on the supporting device from accidentally sliding so that a part of his or her body might perhaps be trapped between an inflated section and, for example, a base section in the process of being inflated, it is preferred that the supporting device be equipped with a sheet of cloth which on one side is fastened to the upper layer of the base and on the other side to the wall at about half way up the height thereof."}
{"text": "1. Technical Field\nThis disclosure relates generally to the field of optics, and in particular but not exclusively, relates to display screens.\n2. Background Art\nConventional rear projection screens use a thin diffusing layer in the form of a matte screen. An image is projected onto the back of the diffusing layer, where it is diffused and scattered into the viewing environment. The diffusing layer provides an image surface and its diffusive nature serves to increase the viewing angles from which the image can be observed. In order to provide good image quality, it is desirable for these projection screens to provide high image contrast and sharp image quality.\nConventional rear projection screens often suffer from poor image contrast in well-lit rooms, at least in part, from the scattering of ambient light back into the viewing environment. This backscattered ambient light deleteriously affects the black levels, and therefore contrast, of the rear projection screen. Furthermore, conventional rear projection screens often suffer from reduced image sharpness or image blur due to stray light in the projection path preceding the image plane of the projection screen. Stray light (i.e., display light with sufficiently large non-normal angles) emitted from the display source at a given display pixel location should be prevented from entering the optical path reserved for an adjacent display pixel prior to the display light reaching the image plane."}
{"text": "Semipermeable membranes for fluid separations are well known in the art. These membranes typically are either dense or asymmetric membranes that include a dense separating layer.\nThe semipermable membranes of the prior art may be polymers of aromatic imides, aromatic esters, and aromatic amides. Membranes of the prior art such as polyester membranes can be derived from bisphenol A and mixtures of iso- and tere-phthaloyl chlorides. Aromatic polyesters prepared from iso/tere-phthalic acids or their halogen derivatives, with substituted biphenols also are known. These polymers, however, are subject to specific structural requirements.\nThe polyester membrane compositions of the prior art, although useful as gas separating membranes, not only suffer from the disadvantages of having to satisfy specific structural constraints, but also are difficult to fabricate into configurations such as hollow fiber membranes since these compositions tend to be soluble in relatively few solvents. Moreover, the membranes of the prior art tend to have relatively low flux.\nA need therefore exists for fluid separation membranes that avoid the fabrication and solubility problems of the prior art but also provide improved gas separation properties."}
{"text": "1. Field of Invention\nThis invention relates to modifying an appearance of an image formed by an image forming device. More particularly, this invention is directed to a method and apparatus for detecting embedded pixels of an image and altering the level of such detected embedded pixels when forming a hard copy of the image.\n2. Description of Related Art\nU.S. Pat. Nos. 5,666,470 to Parker and 5,687,297 to Coonan et al., each co-owned by the assignee of this application and each incorporated by reference herein in its entirety, are directed to image forming systems that allow the appearance of a printed copy of an image to be tuned. Particularly, the various definitions and image appearance tuning systems and methods set forth in the incorporated '470 and '297 patents are specifically incorporated herein.\nIn particular, each of these patents is directed to an image forming device that forms an image based on a bitmap, i.e., a one-bit-depth image data map, or a bytemap i.e., a multi-bit-depth image data map. The appearance tuning systems of these patents apply morphological filters to the bitmap or bytemap to determine which pixels of the bitmap or bytemap require appearance tuning. In particular, the appearance tuning systems disclosed in these patents apply predefined templates to the bitmap or bytemap and scan the templates on a pixel-by-pixel basis over the bitmap or bytemap. Each template defines a different pattern of image data values for a center pixel of interest and one or more adjacent, or \"neighboring,\" pixels of the bitmap or bytemap. If the center pixel of interest and its neighboring pixels do not match any of the predetermined templates, the appearance of the center pixel of interest is not altered. Otherwise, depending on which of the predetermined templates is matched by the center pixel of interest and its neighboring pixels, the appearance of the center pixel of interest is altered to tune the overall appearance of the image to be formed.\nIn particular, the '470 patent is directed to dilation or erosion of the image data value of the pixel of interest to smooth the transition from neighboring black pixels to neighboring white pixels. In particular, dilation refers to increasing the density of the pixel of interest, i.e., making the pixel of interest darker. In contrast, erosion refers to decreasing the density of the pixel of interest, i.e., making the pixel of interest lighter. The '297 patent also modifies the image data value of the pixel of interest to increase or decrease its image data value based on the image data values of its neighboring pixels. In general, the '297 patent differs from the '470 patent primarily in the set of predetermined templates that are used in the template matching or morphological filtering operation."}
{"text": "Calmodulin is an ubiquitous Ca.sup.2+ -binding regulatory protein that mediates the action of Ca.sup.2+ in a diverse array of biological events including cell division, microtubular depolymerization and polymerization, cyclic nucleotide metabolism, Ca.sup.2+ transport, and protein phosphorylation .sup.6,8,31,38,43,72. The involvement of calmodulin in various nerve cell processes, together with its abundant distribution in brain tissue.sup.75, suggests that calmodulin plays an important role in the growth and differentiation of nerve cells.\nThe notion that calmodulin is a regulator of cell growth was initially suggested by studies showing that calmodulin antagonists inhibited cell division in vitro by blocking cell-cycle progression at the G1/S boundary and during G2/M.sup.6,15,60. Consistent with such results are the findings that calmodulin levels are altered during the various phases of the cell cycle.sup.7,42,60 and that calmodulin increases DNA synthesis in a variety of cells .sup.2,24,41,42. Using eukaryotic expression vectors to produce sense or antisense calmodulin RNA, it was demonstrated that increases or decreases in the intracellular levels of calmodulin result in increases and decreases, respectively, in the rates of proliferation of mouse C127 cells.sup.53,54. The role of calmodulin in cell growth is also supported by the findings of elevated levels of calmodulin in certain transformed cell types.sup.12,32,66,69.\nAlthough there is evidence for the participation of calmodulin in cell differentiation, its role has not been precisely defined. Treatment of leukemic HL-60 cells with calmodulin antagonists resulted in an increase in cellular differentiation.sup.67 . On the other hand, the levels of calmodulin were found to increase during nerve growth factor (NGF) - induced differentiation of PC12 cells.sup.4.\nStructural studies have revealed that the amino acid sequence of calmodulin is remarkably conserved among species.sup.50. In the rat genome, calmodulin is encoded by three genes expressing five transcripts, each of which encodes an identical protein, the sequence variations in the coding regions being within the limits of codon degeneracy .sup.47,48,50,63. These calmodulin mRNAs are heterogeneously distributed both in the newborns.sup.11,70 and adult rat brain.sup.18,56 and have been shown to change during ontogeny in a pattern which is paralleled by similar changes in the biological activity of calmodulin.sup.11, indicating that calmodulin is involved in neuronal development.\nUsing the PC12 clonal cell line.sup.21 as an in vitro model for investigating neuronal development, it was demonstrated that PC12 cells express, in varying levels, all five calmodulin gene transcripts.sup.4. Differentiation of PC12 cells induced by NGF or dibutyryl cyclic AMP was accompanied by characteristic and distinct changes in the levels of the various calmodulin mRNAs, demonstrating that the calmodulin genes can be differentially regulated.sup.34. Other studies showed that the multiple calmodulin transcripts have distinct subcellular distributions, some being found both in the cell bodies and neurites and others found only in the cell bodies.sup.74. Since the cellular location of calmodulin will determine its functional activity, these results suggested that the regulation of the levels of particular calmodulin mRNAs in different subcellular compartments may, in turn, be one means for controlling the different functional activities for calmodulin in nerve cells. However, it has not yet been determined if calmodulin gene expression plays a causal role in differentiation."}
{"text": "The present invention relates to diffusion coatings for components exposed to oxidizing environments (i.e., the hostile thermal environment of a gas turbine engine). In some of the more specific embodiments, this invention teaches a process for selectively stripping the additive sublayer of diffusion aluminide coatings from metal substrates.\nIn order to increase the efficiency of gas turbine engines, higher operating temperatures are continuously sought. For this reason, the high temperature durability of the engine components must correspondingly increase. With the formulation of superalloys, such as nickel-base and cobalt-base, significant advances in high-temperature capabilities are being achieved. Consequently, in the absence of a protective coating, sensitive superalloy components (e.g., the turbine and combustor) typically will not endure long service exposures, without wear or damage. One such coating is referred to as an environmental coating, i.e., a coating that is resistant to oxidation and hot corrosion. This coating is typically formed by a diffusion process, e.g., using a pack cementation-type procedure, and usually contains aluminum.\nThe diffusion process generally entails reacting a surface component with an aluminum-containing gas composition. This forms two distinct sublayers, the outermost of which is referred to as the additive sublayer, and the innermost of which is a diffusion sublayer. The additive sublayer contains the environmentally-resistant intermetallic, represented by MAI; where M is iron, nickel or cobalt, depending on the substrate material. The MAI intermetallic is often the result of the diffusion of deposited aluminum into the substrate, and a general, outward diffusion of iron, nickel or cobalt from the substrate. During high temperature exposure in air, the MAI intermetallic forms a protective aluminum oxide (alumina) scale that inhibits oxidation of the coating and the underlying substrate. The chemistry of the additive sublayer can be modified by the presence of additional elements, such as chromium, silicon, platinum, rhodium, hafnium, yttrium and zirconium. As a result of changes in elemental solubility (in the local regions of the substrate and gradient), the diffusion sublayer is thus formed. Due to reactivity, the diffusion sublayer contains various intermetallic and metastable phasesxe2x80x94products of all alloying elements from the substrate and coating.\nThough significant advances have been made with environmental coating materials (and processes for forming such), they still require repair necessitated by erosion and thermal degradation. Moreover, the component on which the coating is formed must often be refurbished. The current state-of-the-art repair method is to completely remove a diffusion aluminide coating, I.e., a xe2x80x9cfull stripxe2x80x9d of both the diffusion sublayer and the additive sublayer.\nAn example of a full stripping process is shown in FIGS. 1 and 2. FIG. 1 is a photomicrograph of a platinum aluminide diffusion coating applied over a superalloy substrate. In this figure, region A is a portion of a substrate, while region B generally represents the diffusion sublayer of a platinum aluminide diffusion coating. Region C is the additive sublayer of the diffusion coating. This additive sublayer serves as a protective layer over the substrate, e.g., a substrate in the form of a turbine wall. As shown in FIG. 1, diffusion sublayer B is incorporated into the upper region of the substrate. (Dotted line E represents the original surface of the substrate).\nFIG. 2 is a photomicrograph of the coated substrate of FIG. 1, after a current, state-of-the art, full stripping technique. Both the additive sublayer and the diffusion sublayer have been removed. In this figure, region D is the remaining portion of the substrate. The original surface of the substrate is again indicated by dotted line E.\nAlthough functional, this full stripping process relies on lengthy exposure to stripping chemicals at elevated temperatures, which causes complete removal of the additive and diffusion sublayers. The process may also significantly attack the underlying metallic substrate, resulting in alloy depletion and intergranular/interdendritic attack. This attack is most prevalent when a component being stripped has regions with different coating thicknesses, or has uncoated surface regions (e.g., the dovetail of a turbine blade). One unfortunate result may be the need for continuous, uneconomical coating replacements.\nAccordingly, it is sometimes undesirable to remove a significant portion of the coating as shown in FIG. 2. For gas turbine blades and vane airfoils, removal of the diffusion sublayer can lead to excessively thinned walls, and drastically altered airflow characteristics.\nIt is apparent that a new process for removing coatings from substrates (i.e., metal substrates) should be welcome in the art. The process must be capable of removing only the additive sublayer of the coating, while not substantially affecting the diffusion sublayer.\nThe present invention generally provides a method of removing an additive sublayer of a diffusion aluminide coating from a metallic substrate. The process described herein is effective In selectively removing only the additive sublayer of a diffusion-aluminide coating. Such a procedure does not attack the underlying diffusion sublayer, including the substrate (see FIGS. 3 and 4, discussed below). In brief, FIG. 3 is a photomicrograph of a coated substrate similar to FIG. 1, prior to the stripping treatment according to the present invention. FIG. 4 depicts the same substrate, after treatment.\nThe present method comprises the step of contacting the diffusion aluminide coating with an aqueous composition comprising at least one acid having the formula HXZrF6, or precursors to said acid, wherein x is 1-6. The acid is typically present at a concentration in the range of about 0.05 M to about 5 M, where M represents molarity. (Molarity can be readily translated into weight or volume percentages, for ease in preparing the solutions). As described below, these compounds may sometimes be formed in situ.\nAn additional embodiment of the invention is directed to an aqueous composition used for the partial stripping of diffusion aluminide coatings. The substrate is a metallic material, and is often in the form of a superalloy component. As used herein, xe2x80x9cpartial strippingxe2x80x9d of the coating refers to the selective removal of substantially all of the additive sublayer, while removing little or none of the underlying diffusion sublayer. The composition comprises an acid having the formula Hx ZrF6, as mentioned above. The acid is usually present in the composition at levels described hereinafter.\nFurther details regarding the various features of this invention are found in the remainder of the specification."}
{"text": "A pressure relief valve (PRV) is a safety device which is designed to divert pressure in the event that a pressurized system of some kind encounters a problem which leads to a dangerous rise in pressure. In the event that pressure in the system becomes too high, instead of blowing out the entire system or damaging system components, the pressurized liquid or gas will vent from the pressure relief valve, bringing the pressure back down and preventing a serious incident. Pressure relief valves can be set to open at a specific pressure.\nIn one application, relief valves may be incorporated into a PRV assembly to be used on a fuel line of an engine system, to release fluid and maintain pressure in the system based on a preset pressure value."}
{"text": "The present invention relates generally to heterojunction bipolar transistors (HBTs), and, more particularly, to HBTs constructed to have reduced thermal resistance to permit rapid dissipation of heat.\nIn HBTs operating at high power, the performance of the device is limited by the amount of power that can be dissipated in the device. Typically, a high power HBT will generate a large amount of heat which must be dissipated quickly to permit proper operation.\nRecently, metallic thermal shunts have been developed to permit dissipation of heat from the collector of the HBT device. Typically, in any HBT, and especially in InP based HBTs, most of the power dissipation and heat generation are located in the upper region of the collector layer near the base layer, where there is a large voltage drop and a large current. FIG. 1 shows an HBT with a metallic thermal shunt designed in accordance with known principles to dissipate this heat.\nReferring to FIG. 1, an HBT 110 is shown. This HBT is constructed with a substrate 114 formed of InP, a sub-collector 116 formed on the substrate, a collector 118 formed on the sub-collector, a base 120 formed on the collector and an emitter 122 formed on the base. An emitter metal 124 is provided on the emitter 122 (with collector metal being provided on the sub-collector 116 and base metal provided on the base 120).\nAs can be seen in FIG. 1, the majority of the heat generated during operation of the HBT 110 is located in a region 112 near the base-collector junction. In order to dissipate this heat to the substrate, a metallic thermal shunt 126 is coupled between the emitter metal 124 and the substrate 114. Thus, the heat generated in the collector 118 will dissipate through the base 120, the emitter 122, the emitter metal 124 and the metallic thermal shunt 126 to the substrate 114. It is noted that this arrangement is also taught in U.S. Pat. No. 5,734,193 to B. Bayraktaroglu et al.\nFrom their studies of the thermal shunt technique described above, the inventors have noted certain disadvantages. The first of these is that the heat dissipation is limited by the fact that the heat must go through the base and emitter layers (and an emitter cap layer if one is used), and these layers generally have relatively low thermal conductivity, especially in InP-based HBTs. For example, the material used in the base of InP-based HBTs is the ternary InGaAs that has a very low thermal conductivity. Furthermore, in InP-based HBTs, the emitter material is typically the ternary semiconductor InAlAs, which also has a very poor thermal conductivity. The thermal shunt approach would provide limited improvement for these type of devices. Also, InGaAs is generally used as a cap layer on top of the emitter to improve the ohmic contact resistance to the emitter, which adds another layer with low thermal conductivity in the path for heat dissipation.\nThe metallic thermal shunt technique was initially developed for AlGaAs/GaAsxe2x80x94based HBTS. In these HBTs, the emitter layer consists of the ternary semiconductor AlGaAs, which has a poor thermal conductivity (see FIG. 4). Therefore, the thermal shunt technique also provides limited improvement in this instance.\nA second disadvantage of the thermal shunt approach is the fact that it requires an extra processing step to construct the thermal shunt after the device itself is completed. Also, additional space is required for the shunt, thereby increasing the overall size of the device.\nIt is an object of the present invention to provide an HBT with improved heat dissipation.\nIt is a further object of the present invention to provide an HBT in which heat is dissipated downward to the substrate directly without the need for a thermal shunt.\nIt is still a further object of the present invention to provide an HBT with reduced thermal resistance between the heat generated near the collector-base junction and the substrate.\nYet another object of the present invention is to provide a method of fabricating an HBT to improve heat dissipation downward to the substrate.\nTo achieve these and other objects, an HBT is provided with a collector and a sub-collector which are each comprised of InP and located relative to one another so that heat generated in the collector during operation of the HBT dissipates downward from the collector through the sub-collector into the substrate."}
{"text": "1. Field\nThis disclosure is generally related to analyzing heterogeneous events. More specifically, this disclosure is related to a method and system for analyzing a large set of event data by introducing a probability model to cluster heterogeneous events.\n2. Related Art\nFor many applications, it is useful to analyze heterogeneous, information-rich events. Heterogeneous events are events that may vary by different factors, including event type, descriptors, location and time. For example, one type of heterogeneous event can be found in military applications. The military may monitor field operations that produces events such as meetings between people of interest, field reports filed by personnel, images and sounds recorded by equipment deployed in locations of interest, and improvised explosive device (IED) explosions.\nDepending on context, analysts may classify events as shallow or deep. Shallow events are those for which relatively little information is available beyond event type, location, and time. Deep events are those for which a rich set of information is available, such as a long field report or a video sequence capturing the event.\nSystems for analyzing event data may collect homogenous or heterogeneous event data. When events are homogenous, all events are of the same type (e.g., observing a pine tree of a particular species) and are characterized by the same set of descriptors (e.g. the girth, height, and age of the tree). Another example of a homogenous event is a “check-in” event where certain software applications may produce events when users check into a venue at a certain time and location.\nWhen the events are heterogeneous, multiple event types are present (e.g. meetings, patrols, and IED explosions), and each event is characterized by a potentially different set of descriptors. For example, an IED detonation can be characterized by descriptors such as power and materials used. These descriptors are inapplicable to other events such as meetings between people, which is characterized by a different set of descriptors (e.g., the set of people involved and the meeting duration). Modeling heterogeneous events is particularly important when there are interactions between events (e.g. meetings between suspected terrorists may precede planting an IED).\nCurrent approaches for analyzing events cannot scale to very large data sets. Furthermore, such approaches typically only deal with homogenous events where arguments are of the same type and described by the same set of descriptors. To process data sets containing heterogeneous events, one can split the heterogeneous events into several data sets, each restricted to a particular type of event. This approach, however, does not allow for modeling interactions between events.\nSome researchers have proposed techniques for analyzing interactions between heterogeneous events. Such proposals typically involve manually constructing detectors for specific activities of interest. For example, some researchers have proposed detectors for intrusion detection. However, such manually constructed detectors may not be used for exploratory analysis. Furthermore, each new detection task requires manually developing a dedicated detector. This is a complex, time-consuming and error-prone procedure."}
{"text": "Gas streams often carry particulate material therein. In many instances it is desirable to remove some or all of the particulate material from the gas flow stream. For example, air intake streams to engines for motorized vehicles or power generation equipment often include particulate material therein. The particulate material, should it reach the internal workings of the mechanisms involved, can cause substantial damage. It is therefore preferred, for such systems, to remove the particulate material from the gas flow upstream of the engine or other equipment involved. A variety of air cleaner arrangements have been developed for particulate removal.\nThere has been a general trend for the utilization of air cleaner arrangements that utilize, as a media pack, z-filter media constructions. In general z-filter media constructions can be characterized as comprising a fluted sheet secured to a facing sheet, formed into a media pack configuration. Examples of z-filter arrangements are described in PCT Publication WO 97/40918, published Nov. 6, 1997; U.S. Pat. Nos. 6,190,432 and 6,350,291; PCT application US 04/07927, filed Mar. 17, 2004; U.S. Provisional application 60/532,783, filed Dec. 22, 2003; PCT Publication 03/095068, published Nov. 20, 2003; PCT publication WO 04/007054, published Jan. 22, 2004; PCT publication WO 03/084641, published Oct. 16, 2003; and, U.S. Provisional Application 60/543,804, filed Feb. 11, 2004; the complete disclosures of each of these cited references being incorporated herein by reference.\nIn general, improvements have been sought."}
{"text": "1. Field of the Invention\nThe present invention relates to substrate processing, and more particularly to improving the substrate processing using optimized Multi-Layer/Multi-Input/Multi-Output (MLMIMO) models, procedures, and subsystems.\n2. Description of the Related Art\nEtch process behavior is inherently non-linear and interacting step-to-step (layers) or as process stacks are compiled (etch/cvd/implant). With the knowledge of the process interactions based on physical modeling of Tokyo Electron Limited (TEL) chambers and base processes and imperial data and measurements from process refinement and tuning the control of Critical Dimension (CD), Sidewall Angle (SWA), depths, film thicknesses, over etching, undercuts, surface cleaning and damage control can be recursively calculated and optimized using multi-input multi-output non-linear models. Current low cost products use a bulk silicon technology. As the transistor continues to shrink, the impact of the channel depth is becoming critical (ultra-shallow source/drain extensions). As the SOI film shrinks, smaller variations in the gate and/or spacer thickness and thickness of the SOI [define SOI] film can affect the transistor's performance. When etch procedures are not controlled, the removal of the material near the gate affects the electrical performance.\nCurrent high performance microprocessors use PD SOI (partially depleted Silicon-on-Insulator film—giving a threshold voltage 0.2 volts. PD SOI films are around 50 nm while the gate and/or spacer reduction amount can be a large percentage (10%) of the total gate and/or spacer thickness. Future generations of SOI films are called FD SOI (fully depleted giving a threshold voltage 0.08 volts and a thickness of ˜25 nm). Currently theses films are not in production due to limitations in thickness control uniformity and defects. Channel mobility degrades with decreasing SOI thickness. With thinner SOI film, the control of the gate and/or spacer sidewall thickness becomes more critical."}
{"text": "The present invention relates generally to structural or building components and in particular to mastic compositions, methods of forming same, and composite structural panels utilizing mastic compositions as well as methods and apparatus for forming such panels.\nIn the past, various mastic compositions, methods of forming them, and composite structural panels utilizing mastic compositions as well as methods and apparatus for forming composite structural panels have been disclosed, but while each had certain advantageous features, none contained all or even most of the features desirable for forming or use of the composite panels as a structural or building component.\nU.S. Pat. No. 3,021,291 discloses mixing expandable polystyrene beads with a cementitious slurry, using a gas to form a cellular material with voids therein, and, after setting or curing, heating to expand the beads which are intended to fill voids. U.S. Pat. Nos. 3,630,820 and 3,697,366 each disclose that closed-cell expanded polystyrene beads, when mixed in a cementitious slurry, rise or float to the top of the slurry so as to make it impossible to attain a homogeneous mixture. In these disclosures, expanded polystyrene beads are shredded to open the cells of the normally closed-cell beads thereby to permit them to absorb slurry into their shredded open cells which overcomes the flotation problem and ensures a homogeneous mix of shredded open-cell beads with the slurry. U.S. Pat. No. 3,598,672 describes mixing expanded polystyrene beads with an epoxy resin, admixing glass fibers which adhere to the surfaces of the beads, subjecting this mixture to pressure to effect setting thereof and then separating the beads coated with glass fibers for use in forming other products.\nU.S. Pat. No. 2,633,439 discloses a composite structural panel with a core of spiral wood shavings and gypsum plaster facing layers or walls with the edges of the shavings embedded in the opposite wall surfaces, but each of the spiral wood shavings has a central aperture to provide free passage of air between cells and to the exterior of the panel. U.S. Pat. No. 2,892,339 shows a honeycomb core structural panel formed in its entirety of molded gypsum material. One face and the honeycomb structure are integrally formed, then cured, and then pressed into the other face to form the panel. One notch is provided in one wall edge of each individual honeycomb cell to permit the circulation of air throughout the core of the assembled honeycomb structure, in which wire mesh and/or fibers may be used for reinforcement. U.S. Pat. No. 3,249,659 discloses the application of a thin layer of foam-forming polyurethane resin to opposite surfaces of two webs of paper facing material which are brought into contact with opposite faces of a honeycomb core and the assembly then heated to effect foaming of the resin. Holes are provided in the honeycomb walls interconnecting each individual honeycomb cell to permit air circulation throughout the entire honeycomb core."}
{"text": "1. Field of the Invention\nThe present invention generally relates to ricin toxin. In particular, the present invention relates to ricin vaccines, compositions and therapeutics as well as methods of making and using thereof.\n2. Description of the Related Art\nRicin is a very toxic protein obtained from the castor bean, Ricinus communis, Euphorbiaceae. Ricin is a heterodimer comprising an A chain and a B chain joined by a disulfide bond. Ricin A chain (RTA) is an N-glycosidase enzyme that irreversibly damages a specific adenine base from 28S rRNA. Once the rRNA has been damaged, the cell cannot make protein and will inevitably die (cytotoxicity). As RTA exhibits this type of destructive catalytic activity, RTA is commonly referred to as a type II ribosome inactivating protein (RIP). See Lord, et al. (1991) Semin. Cell Biol. 2(1):15–22. RTA has been coupled with a targeting moiety to selectively destroy target cells such as tumor cells. See U.S. Pat. Nos. 4,980,457; 4,962,188; and 4,689,401; see also Vitetta et al. (1993) Trends Pharmacol. Sci. 14:148–154 and Ghetie & Vitetta (1994) Cancer Drug Delivery 2:191–198.\nThe toxic consequences of ricin are due to the biological activity of RTA. Ricin B chain (RTB) binds the toxin to cell surface receptors and then RTA is transferred inside the cell where inhibition of ribosome activity occurs. The human lethal dose of ricin toxin is about 1 μg/kg. As highly purified ricin is readily available using methods known in the art, the use of ricin toxin in biological warfare and terrorism is highly possible and probable.\nRicin toxin (RT) or Ricin communis agglutinin II (RCA 60), a glycoprotein produced by the castor bean plant, Ricin communis, is composed of two subunits, about a 30 kDa enzymatically active A subunit (RTA) and about a 32 kDa B subunit (RTB). See Lord & Roberts (1994) Faseb J. 8(2):201–208. The B-chain mediates receptor binding of the toxin to eukaryotic cells via its high affinity for galactose. See Alami & Taupiac (1997) Cell Biol. Int. 21(3): 145–150. Once internalized within the cell, the A chain causes catalytic depurination of the 28S ribosomal RNA that results in inhibition of protein synthesis. See Chen & Link (1998) Biochemistry 37(33):11605–11613. Ricin is highly toxic and can cause death when given in sufficient quantities by either systemic or inhalational routes of exposure. See Wilhelmsen & Pitt (1996) Vet. Pathol. 33(3):296–302.\nRicin is a Category B Agent on the Centers for Disease Control (CDC) Select Agent List and thus there is a strong interest in developing diagnostic tests for toxin identification in clinical and environmental samples. See Thomas, M. (2002) “Possession, use, and transfer of select agents and toxins; interim final rule.” Federal Register 240(67). In addition, because there is no vaccine for ricin and no therapeutic agents available for treatment, there is a serious need to develop prophylactic and therapeutic countermeasures for ricin intoxication.\nDevelopment of antibodies recognizing determinants on the ricin molecule may be able to address several of these needs. Not only can antibodies be used for diagnostic reagents, but they can also neutralize the toxin by either preventing binding to cells or inhibiting enzymatic activity. There is evidence to suggest that antibodies can protect against ricin intoxication as animals were protected from lethality by administration of polyclonal antibody prior to exposure to ricin. See Hewetson & Rivera (1993) Vaccine 11(7):743–746; and Houston (1982) J. Clin. Toxicol. 19(4):385–9. Anti-ricin IgG has also been shown to protect against inhalational challenge in animals, demonstrating the feasibility of using antibody to protect against this route of exposure as well. See Griffiths & Lindsay (1995) Hum. Exp. Toxicol. 14(2):155–164; and Poli & Rivera (1996) Toxicon. 34(9):1037–1044.\nAlthough polyclonal antibody can be used for these purposes, monoclonal antibodies offer several potential advantages, including consistency and reproducibility of product and the ability to humanize the antibody molecule to reduce adverse reactions caused such as serum sickness when animal antibodies are used therapeutically. Several monoclonal antibodies (Mab) previously developed have been shown to confer protection against ricin intoxication in vitro. See Colombatti & Johnson (1987) J. Immunol. 138(10):3339–33344; Colombatti & Pezzini (1986) Hybridoma 5(1):9–19; and Columbatti (1997) Personal communication to M. Dertzbaugh. However these Mabs were lost several years ago and to date only one Mab that still exists which has been shown to protect against ricin intoxication in vivo and this Mab is directed towards the A chain of the holotoxin. See Lemley & Amanatides (1994) Hybridoma 13(5):417–421.\nThus, a need exists for Mabs against ricin toxin."}
{"text": "In recent years, in order to control the occurrence of smoke, a non-burning type flavor inhaler (for example, a smokeless tobacco) that enables inhaling of flavors without burning has been provided (for example, Patent Literature 1). Specifically, the non-burning type flavor inhaler includes a heat source having a cylindrical shape, and a flavor generating source that has a columnar shape and that is provided on the inner side of the heat source. Alternatively, the non-burning type flavor inhaler includes a heat source having a columnar shape, and a flavor generating source that has a cylindrical shape and that is provided on the outer side of the heat source.\nHowever, in the above-described non-burning type flavor inhaler, it is necessary to process the heat source or the flavor generating source in a cylindrical shape, and the manufacture of the heat source or the flavor generating source is complicated. Furthermore, in a case of the above-described cylindrical non-burning type flavor inhaler, in order to realize a sufficient heat generation continuity, the heat source is required to have a predetermined volume, and on the other hand, in order to realize a sufficient supply of the flavor, the flavor generating source is required to have a predetermined volume, which makes it difficult to realize a flavor inhaler that is thin and has excellent portability and handling."}
{"text": "1. Field of the Invention\nThe present invention relates to a powder conveying apparatus that conveys powder contained in a powder container to a discharging unit, an image forming apparatus that includes the powder conveying apparatus, and the powder container.\n2. Description of the Related Art\nElectrophotographic image forming apparatuses, such as copiers, printers, facsimile machines, or multifunction peripherals having functions of copiers, printers and facsimile machines, generally form images by causing developing devices to form toner images with developer called toner or carrier. In such image forming apparatuses, toner is consumed through image formation; therefore, in general, toner cartridges containing toner are attached to the image forming apparatuses and when the toner cartridges become empty of the toner, the toner cartridges are replaced with new ones in order to replenish new toner.\nIn a toner supply system using the cartridges as described above, there is a user's demand to use up all toner in the cartridges in order to reduce running costs. Therefore, various methods have been employed, such as a method in which what is called a screw bottle that is a cylindrical container provided with an inner-mounted spiral protrusion is used as a toner cartridge and toner is gradually conveyed to a discharging unit by rotating the container, or a method in which a screw typically called an auger is provided inside the container and toner is conveyed to a discharging unit by rotating the screw.\nIn the conveying system using the auger, it is needed to arrange and rotate the screw inside the container, so that the configuration becomes complicated. Furthermore, in this conveying system, because a stack of toner is forcibly conveyed by the auger, load is applied to the toner and the toner may be aggregated or deteriorated. Moreover, arranging the screw inside the container that is a replaceable part leads to increase in costs of consumables, so that environmental loads increase because of resource consumption.\nOn the other hand, in the conveying system using the screw bottle, it is not needed to arrange a screw inside the container. Therefore, the configuration becomes simple. However, in this conveying system, because the container itself is rotated when it is used, the container usually has a shape of a cylinder with an outlet arranged on one side surface of the body thereof (a shape like a bottle being laid down). Therefore, the container is disadvantageous in that the capacity for housing toner becomes smaller than a container in a rectangular-solid shape or the container may be too slippery for a person to hold when the container is replaced.\nBesides, the container provided with the inner-mounted auger or the screw bottle is constructed of a container called a “hard bottle” that is made of polyethylene terephthalate (PET) or the like. Therefore, there is a problem in recycling used containers. Specifically, while used containers are generally collected from users by manufacturers and then recycled, reused, or incinerated, because it is difficult to squeeze the containers of the hard bottle type to reduce the volume of the containers, the containers become bulky and costs for collection and transportation increase. Furthermore, when the collected containers are refilled with toner for reuse, there is a difficulty in cleaning the collected containers and toner filling efficiency is not good. Therefore, costs for reusing the collected containers also increase.\nAlternatively, as a toner conveying system that does not use the screw bottle and the auger, there is a method in which a container is caused to oscillate (reciprocating movement) by applying shock to the container from outside or by bringing the container into contact with a stopper so that toner can be moved and discharged with the aid of the inertia force thereof (see Japanese Patent Application Laid-open No. 2002-46843, Japanese Patent Application Laid-open No. 2002-268346). In this system, when a large amount of toner is contained, the toner stacked in the container collectively moves by the oscillation, so that a satisfactory toner conveying speed can be assured per reciprocating oscillation. However, as the amount of toner in the container decreases, a stack of the toner collapses and the toner is thinly spread, so that the toner conveying speed per reciprocating oscillation decreases with a decrease in the height of the stack of the toner. Therefore, the conveying speed cannot be maintained. Furthermore, a writing system may be influenced by the oscillation of the container and image distortion may occur. Moreover, it takes a long time to fluidize the toner or the toner may be aggregated (blocked) due to an action similar to tapping that occurs by the oscillation.\nThere is a known developer supply device, in which a nozzle is inserted into a flexible container that houses toner and the toner is sucked out by a pump through the nozzle so as to be supplied to a developing device (see Japanese Patent Application Laid-open No. 2005-91879). In the toner supply system of this type, the volume of the flexible container is automatically reduced as the toner is sucked out by the pump. Accordingly, when the container becomes almost empty of the toner, the container is in a squeezed state. Therefore, costs needed to collect, transport, and reuse the used container can be reduced. Furthermore, because the volume of the container is reduced in accordance with the amount of toner remaining in the container, it is advantageous in that the amount of remaining toner can be recognized by the appearance of the container.\nHowever, in the system in which the toner is sucked out by the pump, a discharge port of the container needs to be arranged with face down so that the toner can be easily sucked out through the discharge port (suction port). If the discharge port is arranged sideways and the container is placed in the horizontal direction, it is impossible to collect the toner to the vicinity of the discharge port by gravity, and the toner may be cross-linked and remain in the container without being discharged. Therefore, the container of this type, in which the toner is sucked out by the pump, cannot be placed in the horizontal direction, and the ways to place the container are largely limited. Specifically, to smoothly discharge the toner from the container, the container needs to be inclined by 50° or more with respect to the horizontal plane when the container is placed. Therefore, in the configuration using the container of this type, it is needed to ensure an installation space for a container that is practically long in a vertical direction. As a result, it is difficult to reduce the size of the entire image forming apparatus in the vertical direction.\nAs a method for discharging toner from a flexible container without using the pump as described above, a method has been proposed in which a convex member is pressed against and moved along a container from the outside of the container such that contained toner is pushed out through a discharge port (see Japanese Patent Application Laid-open No. H11-143195). With this method, it is possible to discharge toner even when the container is placed so as to extend in the horizontal direction.\nHowever, in the configuration in which the toner is discharged by pushing the container by the convex member as described above, if the toner is packed due to the pressing action of the convex member, toner discharging may be inhibited. In the worse case, blocking may occur due to aggregation of the toner and the toner cannot be discharged from the container.\nJapanese Patent Application Laid-open No. 2006-258926 discloses an image forming apparatus that includes a toner discharge port formed on a bottom plate of a toner cartridge; a shutter that closes the discharge port; and a shutter open-close mechanism that opens and closes the shutter. The shutter open-close mechanism opens the discharge port in accordance with an operation of attaching the toner cartridge to an apparatus main body and closes the discharge port in accordance with an operation of pulling the toner cartridge out of the apparatus main body. The shutter open-close mechanism includes engagement concave portions formed on left and right side surfaces of the shutter in the attachment direction; and engagement pins as engaging members that are arranged on the apparatus main body side and that are engaged with and disengaged from the engagement concave portions. The engagement pins engaged with the engagement concave portions are supported by an engagement maintaining means to maintain the engaged state.\nWhen the toner cartridge is inserted and reaches near a predetermined position, the engagement pins arranged on an oscillation tray enter guide grooves arranged on a bottom plate member. At this time, the engagement pins are biased in a direction along the inner wall surfaces of the guide grooves by a spring and move relative to the movement of the toner cartridge. Inclined surfaces that shift the engagement pins outward are formed on the guide grooves, and when the inclined surfaces come into contact with the engagement pins, the engagement pins move outward. Engagement aid protrusions are formed at end positions of the inclined surfaces on the outer wall surfaces of the guide grooves. When the engagement pins pass over the inclined surfaces, the engagement pins receive a moving force from the spring and move inward. At the same time, the engagement pins come into contact with the engagement aid protrusions, so that the engagement pins are reliably engaged with the engagement concave portions of the shutter. With this engagement, the movement of the shutter is stopped while the movement of the toner cartridge is continued, so that the shutter is gradually opened to thereby open the toner discharge port.\nOn the other hand, when the toner cartridge is pulled out, the shutter is gradually closed. When the engagement pins engaged with the engagement concave portions are returned to the positions of the engagement aid protrusions, the shutter completely closes the toner discharge port. After the shutter closes the toner discharge port, the engagement pins come off from the engagement concave portions. Therefore, the discharge port is completely closed when the toner cartridge is pulled out, so that it is possible to prevent toner dispersion or toner leakage at the time of replacement.\nHowever, if the toner cartridge is repeatedly replaced while the image forming apparatus is used over time, the spring is repeatedly expanded and contracted by the engagement and disengagement of the engagement pins with the engagement concave portions along with the shutter open-close operation. Accordingly, elastic fatigue occurs on the spring and the spring may be deteriorated. If the spring is deteriorated, it becomes difficult to bias the engagement pins by the spring with a desired biasing force, so that the engagement pins may not be engaged with and disengaged from the engagement concave portions. As a result, it becomes difficult to successfully open or close the shutter over time.\nAs a conveying system different from the above conveying systems, there is proposed a system in which a deformable container is used and a delivery member is pressed against and moved along the container from outside to discharge contained toner (see Japanese Patent Application Laid-open No. H11-143195 (Japanese Patent No. 3548402)). With this conveying system, it is possible to convey the toner with small stress, prevent aggregation or deterioration of the toner, and prevent occurrence of abnormal images due to large oscillation or shock. Furthermore, because it is possible to reduce the size of the container, when used containers to be recycled are collected and transported from user sites to manufacturers because of replacement of cartridges or bottles, the collection and the transportation can be performed easily at lower costs.\nThere is also a known developer supply device in which a nozzle is inserted into a discharge hole arranged on a flexible container, developer sucked out by a pump through the nozzle is supplied to a developing device, and the volume of the flexible container is automatically reduced in accordance with the supply of the developer (see Japanese Patent Application Laid-open No. 2006-085067).\nSuch a flexible powder container that does not have a conveying member inside thereof is advantageous in that the volume of the container can be reduced at the time of collection. However, to ensure the powder conveyance, it is difficult to place the powder container such that a longitudinal side extends in the horizontal direction. Therefore, it is difficult to place the powder container parallel to the developing device to reduce the size of the entire image forming apparatus. Specifically, if a spiral groove is formed on the flexible powder container and the container is rotated to supply developer, the container is twisted by the rotation and the developer cannot be conveyed. Furthermore, if developer is to be conveyed without arranging a delivery member inside the container, it is difficult to obtain an angle at which the gravity is utilized to move the developer. As a result, the developer may be cross-linked and remains in the powder container without being discharged. As described above, because it is difficult to place the conventional flexible developer container (powder container) in the horizontal direction, the container needs to be inclined toward the discharging unit by an angle (normally, 50° or greater) slightly greater than the repose angle of the powder. Therefore, the container that is practically long in the vertical direction needs to be arranged, so that the shape of the image forming apparatus, the capacity of the container, and the arrangement of the container are largely limited.\nIn the above-mentioned system in which the delivery member is pressed against and moved along the deformable container from outside in order to discharge the contained toner, air in the container is also pushed out by the stack of the toner in the container in some cases, and a greater amount of toner than needed may be discharged by the airflow. To prevent the excessive amount of toner from being discharged by air, an air filter that releases air to the outside may be arranged on a toner discharge path. However, the air filter may be gradually clogged and needs to be replaced periodically."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates to the field of mobile communication technology, and more particularly, to a method of accessing to networks and a mobile communication terminal.\n2. Description of the Related Art\nPublic land mobile network (PLMN) is a network built and operated by the government or operators approved by the government for the purpose of providing land-bound mobile communication services for the public. The network is usually connected to public switched telephone network (PSTN), forming a communication network that covers a whole region or country.\nA network identity of the PLMN is usually just a string of numbers. For example, the network identity of PLMN of China Mobile is 46000, and the network identity of PLMN of China Unicom is 46001. Through the mobile networks provided by each operator, users can enjoy rapid and convenient communication. The number of users carrying communication terminals has been on the rise as people's living standard elevates. When a user inserts a new user identification card into a mobile communication terminal for the first time, the mobile communication terminal (e.g. a cell phone) must search a network to acquire the network identity of the PLMN. However, when the user inserts the new user identification card into the mobile communication terminal for the first time, it usually takes the terminal a long time (sometimes even as long as ten to twenty minutes) search the network. And, only when the terminal successfully found the network can it connect to the PLMN of the current location. It significantly influences users' communication experience."}
{"text": "At present, for the purpose of smoothly supplying electric power to an IC chip, a chip capacitor is often mounted on the surface of a printed circuit board used as a package substrate.\nThe reactance of a wiring from the chip capacitor to the IC chip depends on frequency. Due to this, as IC chip driving frequency increases, the chip capacitor cannot exhibit sufficient effect even if the chip capacitor is mounted on the surface of the printed circuit board. Considering this, the applicant of the present invention proposed a technique, identified as Japanese Patent Application No. 11-248311, for forming a concave portion on a core substrate so as to contain therein a chip capacitor. Techniques for embedding a capacitor in a substrate are disclosed in Japanese Patent Unexamined Application Publication (to be referred to as “Publication” hereinafter) Nos. 6-326472, 7-263619, 10-256429, 11-45955, 11-126978, 11-312868 and the like.\nPublication No. 6-326472 discloses a technique for embedding a capacitor in a resin substrate made of glass epoxy. With this constitution, it is possible to reduce power supply noise and to dispense with a space for mounting the chip capacitor, thereby making an insulating substrate small in size. Publication No. 7-263619 discloses a technique for embedding a capacitor in a substrate made of ceramic, alumina or the like. With this constitution, the capacitor is connected between a power supply layer and a ground layer to thereby shorten wiring length and reduce wiring inductance.\nHowever, according to the Publication Nos. 6-326472 and 7-263619 stated above, the distance from the IC chip to the capacitor cannot be set too short and the wiring inductance cannot be reduced as currently desired in the higher frequency region of the IC chip. In case of a multi-layer buildup wiring board made of resin, in particular, due to the difference in the coefficient of thermal expansion and a capacitor made of ceramic and a core substrate as well as interlayer resin insulating layers made of resin, disconnection occurs between the terminal of the chip capacitor and a via, separation occurs between the chip capacitor and the interlayer resin insulating layers and cracks occur to the interlayer resin insulating layers. Thus, the multi-layer buildup wiring board could not realize high reliability for a long time.\nOn the other hand, according to the invention of Publication No. 11-248311, if the position at which a capacitor is arranged is deviated, the connection between the terminal of the capacitor and a via cannot be accurately established and power cannot be possibly supplied from the capacitor to the IC chip.\nThe present invention has been made to solve the above-stated problems. It is, therefore, an object of the present invention to provide a printed circuit board including therein a capacitor and having enhanced connection reliability and to provide a method of manufacturing the printed circuit board."}
{"text": "In recent years, a variety of user interfaces for performing information processing have been proposed. As one example thereof, a near field wireless communication technology such as NFC (Near Field Communication) is sometimes used. For example, by moving a wireless communication terminal using NFC closer to a predetermined object, an application is activated and information associated with a position to which the terminal is moved closer is displayed.\nPTL 1 describes a personal digital assistance which simplifies Internet access processing and retrieval processing of Web pages. The personal digital assistance described in PTL 1 includes a detection means for detecting a logo mark including a graphic or a character, and stores registration data which registers a Web page address corresponding to the logo mark. When an image inputted from a camera detects the registered logo mark, the personal digital assistance described in PTL 1 accesses the pre-registered Web page address.\nPTL 2 describes an image display device which provides good visibility even in a situation where vibration occurs in the device. The image display device described in PTL 2, vibration applied with respect to the image display device is detected, and a display position of an image to be displayed on a display unit is moved in the opposite direction of the movement of the vibration.\nIn PTL 3, a portable telephone device which decreases a consumption current during a voice call is described. In the device described in PTL 3, the state where the portable telephone device is disposed in the vicinity of user's ear is detected on the basis of the brightness or the outline of image data which a camera unit obtains.\nIn PTL 4, a camera-equipped mobile terminal device which calls a function using an image generated by the camera is described. The camera-equipped mobile terminal device described in PTL 4 detects a feature which an application sets as a processing target, from a preview image obtained by the camera. Regarding all applications which set the detected feature as a processing target, each of them sets a state of receiving selection operation, and information which indicates the function and a method of the selection operation is displayed on a display unit in each of these applications.\nIn PTL 5, a portable telephone which can be easily found at an incoming call in the case of being placed on a desk, a floor, or the like is described. The portable telephone described in PTL 5 detects a state where the terminal is placed on a floor or a desk with a back camera, and changes a notifying method at an incoming call by the detection."}
{"text": "1. Field of the Invention\nThe present invention relates to polyetherimide/poly(arylene sulfide) blends compatibilized with an epoxy compound, and more particularly relates to reinforced polyetherimide/poly(arylene sulfide) compatibilized blends.\n2. Description of the Related Art\nCompositions comprising blends of poly(phenylene sulfide) and polyetherimdes have been disclosed in U.S. Pat. No. 4,455,410 (Giles).\nCompositions comprising a poly(arylene sulfide) resins, an epoxy compound, and optionally another thermoplastic resin, a reinforcing material, a filler and/or a fire retardant have been disclosed in U.S. Pat. No. 4,528,346 (Sugie). This patent further lists numerous possible other thermoplastic resins, and specifically sets out polyesters, polyamides, polyarylates, polycarbonates, poly(phenylene oxide)s, polyimides, polyketone imides, polyamideimides, poly(ether ether ketone)s, polysulfones and elastomers, and sets out in Table 4, Examples 44 and 45 compositions containing poly(phenylene sulfide), polyamideimide and an epoxy compound.\nCompositions comprising glass reinforced poly(arylene sulfide) resins and polyetherimide which additionally contain at least one polyolefin in an amount to substantially improve the heat distortion temperature of the composition have been disclosed in commonly owned application Ser. No. 08/028,061 (Glaser et al) filed Mar. 8, 1993, now U.S. Pat. No. 5,430,102.\nGlass filled poly(arylene sulfide)/polyetherimide blend compositions exhibit numerous desirable properties, but there has been a need to enhance various properties, including, for example, the flash resistance and ductility of the compositions.\nIt is an object of the present invention to provide glass filled poly(arylene sulfide)/polyetherimide blend compositions which exhibit enhanced physical properties, for example, exhibiting enhanced ductility and morphological stability, while substantially retaining other key properties such as heat distortion temperature."}
{"text": "1. Field\nThe present invention pertains generally to communications, and more specifically to a novel and improved method and apparatus for hashing over multiple frequency bands in a communication system.\n2. Background\nCommunication systems and wireless systems in particular, are designed with the objective of efficient allocation of resources among a variety of users. Wireless system designers in particular aim to provide sufficient resources to satisfy the communication needs of its subscribers while minimizing costs. Efficient use of resources requires prompt assignment of mobile stations to specific frequencies.\nIn a wireless communication system employing a Code Division-Multiple Access (CDMA) scheme or Wideband Code Division Multiple Access (WCDMA) each of the subscriber units is assigned code channels at designated time intervals on a time multiplexed basis. A central communication node, such as a Base Station (BS) or Node B, implements the unique carrier frequency or channel code associated with the subscriber to enable exclusive communication with the subscriber. Time Division Multiple Access (TDMA) schemes may also be implemented in landline systems using physical contact relay switching or packet switching. A CDMA system may be designed to support one or more standards such as: (1) the “TIA/EIA/IS-95-B Mobile Station-Base Station Compatibility Standard for Dual-Mode Wideband Spread Spectrum Cellular System” referred to herein as the IS-95 standard; (2) the standard offered by a consortium named “3rd Generation Partnership Project” referred to herein as 3GPP; and embodied in a set of documents including Document Nos. 3G TS 25.211, 3G TS 25.212, 3G TS 25.213, and 3G TS 25.214, 3G TS 25.302, referred to herein as the W-CDMA standard; (3) the standard offered by a consortium named “3rd Generation Partnership Project 2” referred to herein as 3GPP2, and TR-45.5 referred to herein as the cdma2000 standard, formerly called IS-2000 MC, or (4) some other wireless standard.\nCDMA2000 is an improvement on TIA/EIA-95. It provides a significant improvement in voice capacity and extended data capability and is backward compatible with TIA/EIA-95 mobiles. When a mobile station moves between base stations in a CDMA2000 system the mobiles must register and be assigned a frequency for communication. The frequency assignment occurs during a registration process. Registration includes a hashing process to assign a frequency to the mobile station. The mobile must re-register when changing between base stations, with each change forcing a new hash to a new frequency, and in many cases a new frequency band. Hashing is triggered for any change in the frequency distribution or frequency weights. Frequency distribution and weighting is an important consideration for balancing system loading and ensuring efficient system operation. The mobile station also updates the system overhead information each time a hash is performed. This can result in additional and excessive frequency changes, as every frequency change results in system acquisition and reading system overhead information. Unfortunately, while re-acquiring the system pages messages directed to the mobile station may be missed.\nAccordingly, there is a need for a method and apparatus for hashing mobiles over multiple bands while avoiding unnecessary frequency changes."}
{"text": "Memory devices are typically provided as internal storage areas in a computer. The term memory identifies data storage that comes in the form of integrated circuit chips. There are several different types of memory used in modern electronics, one common type is RAM (random-access memory). RAM is characteristically found in use as main memory in a computer environment. RAM refers to read and write memory; that is, you can both write data into RAM and read data from RAM. This is in contrast to ROM (read-only memory), which permits you only to read data. Most RAM is volatile, which means that it requires a steady flow of electricity to maintain its contents.\nComputers almost always contain a small amount of ROM that holds instructions for starting up the computer, typically called a basic input output system (BIOS). Unlike RAM, ROM generally cannot be written to by a user. An EEPROM (electrically erasable programmable read-only memory) is a special type of non-volatile ROM that can be erased and programmed by exposing it to an electrical charge. EEPROM comprise a large number of memory cells having electrically isolated gates (floating gates). Data is stored in the memory cells in the form of charge on the floating gates. Charge is transported to or removed from the floating gates by specialized programming and erase operations.\nYet another type of non-volatile memory is a Flash memory. A Flash memory is a type of EEPROM that can be erased and reprogrammed in blocks instead of one byte at a time. A typical Flash memory comprises a memory array, which includes a large number of memory cells. Each of the memory cells includes a floating gate field-effect transistor capable of holding a charge. The data in a cell is determined by the presence or absence of the charge in the floating gate. The cells are usually grouped into sections called “erase blocks.” The memory cells of a Flash memory array are typically arranged into a “NOR” architecture (each cell directly coupled to a bitline) or a “NAND” architecture (cells coupled into “strings” of cells, such that each cell is coupled indirectly to a bitline and requires activating the other cells of the string for access, but allowing for a higher cell density). Each of the cells within an erase block can be electrically programmed in a random basis by charging the floating gate. The charge can be removed from the floating gate by a block erase operation, wherein all floating gate memory cells in the erase block are erased in a single operation. It is noted that other types of non-volatile memory exist which include, but not limited to, Polymer Memory, Ferroelectric Random Access Memory (FeRAM), Ovionics Unified Memory (OUM), Magnetoresistive Random Access Memory (MRAM), Molecular Memory, Nitride Read Only Memory (NROM), and Carbon Nanotube Memory.\nEach erase block of a Flash memory device contains a series of physical pages that are typically each written to a single row of the Flash memory array and include one or more user data areas and associated control or overhead data areas. The control/overhead data areas contain overhead information for operation of physical row page and the user data area each overhead data space is associated with. Such overhead information typically includes, but is not limited to, erase block management (EBM) data, sector status information, or an error correction code (ECC). ECC's allow the Flash memory and/or an associated memory controller to detect data errors in the user data area and attempt to recover the user data if possible.\nMany of the internal operations of volatile and non-volatile memories require that the memory perform data comparisons. Typically this data comparison is performed in the context of comparing data that has been read from the memory array with the data that was expected to be read in order to find discrepancies. The internal operations requiring data comparison include, but are not limited to, data write, block erasure, and memory testing. A problem with many modern memory devices and arrays is that, because of their increasing storage density levels and increasing size of each physical page/row of data of the array, this data comparison within the memory device has become a time consuming task and can affect the speed of operation, data throughput, and testing time required.\nFor the reasons stated above, and for other reasons stated below which will become apparent to those skilled in the art upon reading and understanding the present specification, there is a need in the art for circuits and methods to allow fast and efficient comparison of large amounts of data within memory devices."}
{"text": "1. Field of the Invention\nThe present invention is broadly concerned with improved apparatus for converting a conventional hot water heater into a simple water storage container, and to permit heating of domestic water supplies externally of the heater tank by making use of an existing boiler forming a part of a space heating system. More particularly, it is concerned with such apparatus which in preferred forms includes a double-pipe fitting which can be secured directly to an existing hot water tank using only the threaded drain valve opening normally provided therein, thereby completely avoiding the necessity of drilling additional holes through the heater tank wall.\n2. Description of the Prior Art\nThe energy consumed by a typical hot water heater in a residence or small building in order to provide adequate supplies of hot domestic water is generally very considerable; indeed, in some cases one-third or more of the total household energy consumption may be devoted to heating of water. Obviously, the dramatic rise in energy costs experienced during the last decade has focused attention upon all energy intensive devices, and hot water tanks are no exception. Modern day hot water tanks are therefore provided with significant insulation and other energy saving features. Nevertheless, energy consumption by these units is still significant.\nOne proposal for reducing energy usage in the context of hot water heating is to convert an existing hot water heater to a simple storage tank, and to provide heating for the water externally of the tank in an indirect heat exchanger coupled to an existing boiler forming a part of an heating system for the residence or building. In other contexts heat derived from a solar collector system or a waste heat source can be employed. From a theoretical standpoint, the proposal has a number of advantages, principally stemming from the fact that but a single heating device (i.e., the heating system boiler) is employed for heating domestic water supplies as well as water used in the space heating system. Typically, boilers of this type are more efficient than hot water heaters, and therefore reductions in energy consumption are realized.\nWhile the foregoing proposal has a number of advantages, from a practical standpoint the task of converting a conventional hot water heater tank for use in the described system can be a formidable one. To give but one example, hot water heater tanks are normally provided with only a single threaded drain opening proximal to the lower end of the tank, for receiving a spigot-type drain valve. Thus, a straightforward conversion of a normal hot water heater tank requires drilling of at least one additional hole through the tank wall, in order to provide the necessary water circulation lines for moving water to and from the external heat exchanger. Obviously, such modifications of an existing hot water heater tank are beyond the skills of the average homeowner, thereby requiring professional assistance and consequent significant expense.\nAccordingly, while the concept of hot water heater tank conversion and external heating of domestic hot water supplies is an attractive one, the practical difficulties associated with implementation of the idea have prevented its widespread adoption. Thus, any apparatus which would facilitate such conversions would be a decided advantage, and would be economically attractive in view of the potential reductions in energy consumption which could be achieved."}
{"text": "In recent years, as semiconductor devices are made highly integrated and highly dense, circuit interconnections have become finer and the number of layers of multi-layer interconnections has been increased. Aiming at achieving multi-layer interconnection while aiming at finer circuitry leads to film coverage of step geometry (step coverage) being lowered in thin film formation as the number of the interconnected layers increases because surface steps increase while following surface irregularities on a lower layer. Therefore, in order to obtain multi-layer interconnection, this step coverage has to be improved to perform a planarization process at an appropriate time. In addition, since finer optical lithography entails shallower depth of focus, it is necessary to subject a surface of a semiconductor device to the planarization process so that surface steps of irregularities on the surface of the semiconductor device fall within the depth of focus. As the circuitry is made finer, a requirement for accuracy with respect to the planarization process has been raised. Not only in a multi-layer interconnecting process but also at a FEOL (Front End Of Line), as a transistor's peripheral structure is complexed, the requirement for accuracy with respect to the planarization process has been raised.\nIn this way, in a manufacturing process of the semiconductor device, a planarization technique for the semiconductor device surface has become important more and more. In this planarization technique, the most important technique is chemical mechanical polishing (CMP). This chemical mechanical polishing is a process in which a polishing apparatus is used to perform polishing by supplying a polishing liquid containing abrasive grains such as of silica (SiO2) onto a polishing surface of a polishing pad or the like and bringing a substrate such as a semiconductor wafer into sliding contact with the polishing surface.\nThis type of polishing apparatus includes a polishing table having a polishing surface formed of a polishing pad, and a substrate holder, called a top ring or a polishing head, for holding the semiconductor wafer. In a case where such a polishing apparatus is used to polish the semiconductor wafer, the semiconductor wafer is held by the substrate holder and the semiconductor wafer is pressed against the polishing surface at a predetermined pressure. At this time, the polishing table and the substrate holder are moved relatively to each other such that the semiconductor wafer is brought into sliding contact with the polishing surface to polish the surface of the semiconductor wafer to a flat and mirror finish.\nIn such a polishing apparatus, if a relative pressing force between the semiconductor wafer and the polishing surface of the polishing pad during polishing is not uniform over the entire surface of the semiconductor wafer, insufficient polishing or excessive polishing would occur depending on the pressing force applied to any portion of the semiconductor wafer. In order to unify the pressing force applied to the semiconductor wafer, a pressure chamber formed of an elastic membrane (membrane) is provided at a lower part of the substrate holder, and, by supplying a fluid such as pressurized air to this pressure chamber, the semiconductor wafer is pressed against the polishing surface of the polishing pad by means of a fluid pressure via the elastic membrane to perform polishing.\nThe substrate holder is provided with a retainer ring surrounding the semiconductor wafer (e.g., see Patent Literature 1), and when polishing the semiconductor wafer, the retainer ring is pressed against the polishing surface at a predetermined pressure so that the semiconductor wafer held by the substrate holder does not get out of the polishing head. Here, a pressing force of the retainer ring is also an adjustment parameter for adjusting a polishing profile of a periphery of the semiconductor wafer.\nAs the pressing force of the retainer ring is lowered, a phenomenon cannot be prevented that the retainer ring on a downstream side of table rotation is uplifted by a horizontal force from the wafer caused by friction between the wafer and the polishing pad and the semiconductor wafer during polishing cannot not be held, and thereby, the semiconductor wafer slides on the polishing pad to get out to the outside (hereinafter, referred to as sipping out) at a certain pressing force of the retainer ring (hereinafter, referred to as retainer ring pressure). In order that the semiconductor wafer does not slip out, the retainer ring pressure needs to be set to be equal to or more than a lower limit of retainer ring pressure (hereinafter, also referred to as RRP (retainer ring pressure) lower limit) at which the semiconductor wafer can be polished without slipping out. However, the RRP lower limit varies depending on a process type or a polishing condition, and thus, disadvantageously is difficult to determine.\nAs for dealing with this problem, there may be considered a method in which polishing is actually performed so as to lower the pressing force of the retainer ring until the semiconductor wafer slips out to measure the RRP lower limit. However, in this method, because the semiconductor wafer actually slips out, expendables such as the membrane or the retainer ring may be broken in some cases. Such a method would require time also. Further, the RRP lower limit varies depending on the process type or the polishing condition, which involves a need to conduct a test for finding the RRP lower limit every time the process type or the polishing condition is changed. However, it is not realistic to conduct a test for finding the RRP lower limit every time the process type or the polishing condition is changed, considering time and effort are taken.\nIt is desired to provide a polishing apparatus, a control method and, a recording medium capable of preventing an polishing object from slipping out without depending on the process type or the polishing condition.\nA polishing apparatus according to one aspect of this technique, a polishing apparatus for polishing a surface to be polished of an polishing object by sliding the surface to be polished and a polishing member relative to each other, comprising: a pressing unit that presses a back surface of the surface to be polished of the polishing object such that the surface to be polished is pressed against the polishing member; a retainer member that is arranged on an outer side of the pressing unit and presses the polishing member; a storage unit that stores information concerning a condition for preventing the polishing object from slipping out, the condition being defined by use of information concerning a pressing force of the retainer member; and a control unit that acquires information concerning a force of friction between the surface to be polished of the polishing object and the polishing member or information concerning the pressing force of the retainer member, and executes control for adapting to the condition for preventing the slipping-out by using the acquired information concerning the force of friction or the acquired information concerning the pressing force of the retainer member.\nBy doing so, the condition for preventing the polishing object from slipping out is not changed even if the process type or the polishing condition is varied, which makes it possible to prevent the polishing object from slipping out without depending on the process type or the polishing condition.\nA control method according to one aspect of this technique, a control method for executing control by way of referencing a storage unit that stores information concerning a condition for preventing an polishing object from slipping out, the condition being defined by use of information concerning a pressing force of a retainer member, the method comprising: a step of acquiring information concerning a force of friction between a surface to be polished of the polishing object and a polishing member, or the information concerning the pressing force of the retainer member; and a step of executing control for adapting to the condition for preventing the slipping-out by using the acquired information concerning the force of friction or the acquired information concerning the pressing force of the retainer member.\nBy doing so, the condition for preventing the polishing object from slipping out is not changed even if the process type or the polishing condition is varied, which makes it possible to prevent the polishing object from slipping out without depending on the process type or the polishing condition.\nA recording medium according to one aspect of this technique, a recording medium storing therein in a non-transitory manner a program for executing control by way of referencing a storage unit that stores information concerning a condition for preventing an polishing object from slipping out, the condition being defined by use of information concerning a pressing force of a retainer member, the program causing a computer to execute: a step of acquiring information concerning a force of friction between a surface to be polished of the polishing object and a polishing member, or the information concerning the pressing force of the retainer member; and a step of executing control for adapting to the condition for preventing the slipping-out by using the acquired information concerning the force of friction or the acquired information concerning the pressing force of the retainer member.\nBy doing so, the condition for preventing the polishing object from slipping out is not changed even if the process type or the polishing condition is varied, which makes it possible to prevent the polishing object from slipping out without depending on the process type or the polishing condition."}
{"text": "A printed circuit board (PCB) is formed by printing a circuit line pattern on an electrical insulating substrate by using a conductive material such as copper (Cu), and refers to a board right before electronic parts are mounted thereon. In other words, the PCB refers to a circuit board in which the mounting positions of the electronic parts are determined, and a circuit pattern connecting the electronic parts to each other is fixedly printed on the surface of a flat plate so that several types of many electronic devices are densely mounted on the flat plate\nMeanwhile, in recent years, a buried pattern substrate having the reduced thickness and planarized surface has been used for the purpose of high performance and miniaturization of electronic parts.\nFIG. 1 is a sectional view showing a typical buried pattern PCB 10.\nAs shown in FIG. 1, the buried pattern PCB 10 includes a buried pattern groove 2 in the surface of an insulating substrate 1 and a circuit pattern 3 formed by filling the buried pattern groove 2 through a plating process.\nThe PCB 10 having the buried circuit pattern 3 can represent very strong adhesive strength with respect to an insulating member due to the formation structure of a base circuit pattern and a contact part, and the pitches of base circuit patterns and contact parts can be uniformly and finely formed."}
{"text": "From DE 10 2012 022 882 A1, a device for controlling a machine used in the graphic arts industry is known, having a control console, said control console having at least one screen for the display of machine information, wherein the control console is equipped with a plurality of sensors for detecting gestures of an operator operating the machine via the control console, wherein the control console is equipped with a support surface, disposed underneath at least one screen, for receiving a sheet-type printing substrate, wherein at least one sensor is situated such that it is capable of recognizing gestures made by the operator with his fingers on the surface of a printing substrate that rests on said support surface, wherein the sensors are preferably cameras, and wherein the control console is preferably equipped with a colorimeter for examining a sheet-type printing substrate resting on the support surface, and the colorimeter can be controlled by means of the sensors for detecting gestures.\nDE 10 2006 048 539 A1 discloses a device for positioning a measuring head above a printing substrate, including a display device and a motor-driven positioning device for the measuring head, wherein the positioning device for the measuring head can be controlled by a computer which is connected to the display device.\nA camera for detecting at least one gesture performed by a person is known from the document JOSHI, Tejashri J.; KUMAR, Shiva; TARAPORE, N. Z.; MOHILE, Vivek: Static Hand Gesture Recognition using an Android Device, in: International Journal of Computer Applications (0975-8887), Volume 120, June 2015, pp. 48-53, available online at http://research.ijcaonline.org/volume120/number21/pxc3904348.pdf, retrieved on Jan. 4, 2017.\nDE 101 20 487 A1 discloses a control console of a printing machine, comprising at least one support surface for receiving a sheet-type printing substrate and a measuring device for measuring a measured variable of the sheet-type printing substrate placed on the support surface, the measuring device being equipped with a measuring bar extending above the sheet-type printing substrate, and the measuring bar of the measuring device being movable above the sheet-type printing substrate that has been placed on the support surface."}
{"text": "This invention relates to certain novel analogs of the naturally occurring prostaglandins. In particular, it relates to novel 15-substituted-.omega.-pentanorprostaglandins and various novel intermediates useful in their preparation.\nThe prostaglandins are C-20 unsaturated fatty acids which exhibit diverse physiological effects. For instance, the prostaglandins of the E and A series are potent vasodilators (Bergstrom, et al., Acta Physiol. Scand. 64:332-33 1965 and Bergstrom, et al., Life Sci. 6:449-455, 1967) and lower systemic arterial blood pressure (vasodepression) on intravenous administration (Weeks and King, Federation Proc. 23:327, 1964; Bergstrom, et al., 1965, op. cit.; Carlson, et al., Acta Med. Scand. 183:423-430, 1968; and Carlson, et al., Acta Physiol. Scand. 75:161-169, 1969). Another well known physiological action for PGE.sub.1 and PGE.sub.2 is as a bronchodilator (Cuthbert, Brit. Med. J. 4:723-726, 1969).\nStill another important physiological role for the natural prostaglandins is in connection with the reproductive cycle. PGE.sub.2 is known to possess the ability to induce labor (Karim, et al., J. Obstet. Gynaec. Brit. Cwlth. 77:200-210, 1970), to induce therapeutic abortion (Bygdeman, et al., Contraception, 4, 293 (1971) and to be useful for control of fertility (Karim, Contraception, 3, 173 (1971)). Patents have been obtained for several prostaglandins of the E and F series as inducers of labor in mammals (Belgian Pat. No. 754,158. and West German Pat. No. 2,034,641), and on PGF.sub.1, F.sub.2, and F.sub.3 for control of the reproductive cycle (South African Pat. No. 69/6089). It has been shown that luteolysis can take place as a result of administration of PGF.sub.2 .alpha.[Labhsetwar, Nature 230 528 (1971)] and hence prostaglandins have utility for fertility control by a process in which smooth muscle stimulation is not necessary.\nStill other known physiological activities for PGE.sub.1 are in the inhibition of gastric acid secretion (Shaw and Ramwell, In: Worcester Symp. on Prostaglandins, New York, Wiley, 1968, p. 55-64) and also of platelet aggregation (Emmons, et al., Brit. Med. J. 2:468-472, 1967).\nIt is now known that such physiological effects will be produced in vivo for only a short period, following the administration of a prostaglandin. A substantial body of evidence indicates that the reason for this rapid cessation of activity is that the natural prostaglandins are quickly and efficiently metabolically deactivated by .beta.-oxidation of the carboxylic acid side-chain and by oxidation of the 15.alpha.-hydroxyl group (Anggard, et al., Acta. Physiol. Scand., 81, 396 (1971) and references cited therein). It has been shown that placing a 15-alkyl group in the prostaglandins has the effect of increasing the duration of action possibly by preventing the oxidation of the C15-hydroxyl [Yankee and Bundy, JACS 94, 3651 (1972)], Kirton and Forbes, Prostaglandins, 1, 319 (1972).\nIt was, of course, considered desirable to create analogs of the prostaglandins which would have physiological activities equivalent to the natural compounds, but in which the selectivity of action and the duration of the activity would be increased. Increased selectivity of action would be expected to alleviate the severe side effects, particularly gastrointestinal side effects, frequently observed following systemic administration of the natural prostaglandins (Lancet, 536, 1971 )."}
{"text": "In machines such as homopolar motors utilizing high armature current at low voltage, high current-carrying capacity is required for making electrical connections between the rotor and the stator. It is well known in the art to utilize liquid metal current collectors for such machines due to the high current-density capacity of such collectors, for example, 3,000 to 10,000 amperes per square inch, by the liquid metal as compared to approximately one-tenth of this current density in the case of conventional solid carbon brushes. Further, solid brushes exhibit an order of magnitude higher voltage drop than that of liquid metal brushes, even at the lower current density. As a result of this higher voltage drop, solid brushes generate more heat than liquid metal brushes and wear both the brushes and the slip rings down rapidly. To provide the high current carrying capacity required by homopolar motors, a large number of solid brushes are necessary. The current is not necessarily shared equally by all of the brushes resulting in unequal heating and wear of the brushes. Moreover, a failure of one brush allowing debris to be transported around the slip ring may result in catastrophic failure of all the brushes.\nIn high current, high speed operation of homopolar machines, it has been common practice to provide a constant supply of liquid metal by pumping it to the stator and rotor current collector sites. This ensures that the current collectors are continuously wetted.\nU.S. Pat. No. 4,628,221 to Young discloses a homopolar motor with pressurized liquid metal contact. The invention of Young uses a rotor having a circular cylindrical shell utilized as a conductor ring, a stator current collector ring of one polarity encircling one edge of the rotor conductor ring and another stator current collector ring of the opposite polarity encircling the other edge of the rotor conductor ring. Liquid metal is utilized within the cylindrical enclosure to provide continuous electrical contact between the stator current collector and the rotor ring.\nU.S. Pat. No. 2,588,466 to Barnes discloses a unipolar or homopolar generator using a sodium/potassium alloy as a liquid brushes. Structure adapted to employ liquid brushes is also disclosed.\nGreat Britain Patent Application No. 2,203,293 to Parson discloses a homopolar generator comprising a rotor having an armature of conductive, ferromagnetic material. The rotor being surrounded by a ferromagnetic stator, field coils being provided on the stator, the field coils being connected to an alternating current supply to provide an alternative magnetic field.\nEuropean Patent Application No. 0,347,089 discloses a homopolar device along the line of the above GB '293 application with the exception that it utilizes first and second annular elements having conductive and non-conductive sectors. Formation of annular electrical currents is thus avoided."}
{"text": "A. Field of Invention\nThis invention pertains to an air bag passenger restraint system, and more particularly to a housing for the components used to inflate the air bag in such a system.\nB. Description of the Prior Art\nActive passenger restraint systems and more particularly, restraint systems with air bags have proved to be very effective in preventing serious injury and even death in motor vehicles involved in crashes. Typically, these systems include an inflator housing containing chemical agents which may be triggered, for example by a crash, to generate a large volume of gas. The housing is arranged and constructed so that the gas is channeled to a bag which inflates. Thus, the bag cushions the body of the passengers.\nEarly gas generator or inflator designs comprised two metal cups which were held together on their circumference to form a pressure chamber. As pressures typically reach 2000 psi within a gas generator, the wall thickness of these cups had to be sufficient to withstand this pressure typically with a safety factor of 2. These cups were typically held together by welding, threads, or by mechanically bending one cup around the other.\nIn an effort to reduce the overall weight of the gas generator by reducing the wall thickness of the cups, means were sought to provide a physical attachment of the two cup halves at some inner diameter. In one prior art configuration, in addition to the peripheral attachment, an inner bolt circle of rivets was used to provide this additional attachment as well as a separate two-piece threaded member that served both to hold the electric squib as well as provide mechanical attachment of the upper and lower half. The outer periphery of this type of inflator was welded.\nAnother configuration utilized a butt welding technique wherein the two cup halves were friction-welded together at both the peripheral diameter as well as an inner diameter. Yet other configurations utilized several interconnecting cup members which were laser or electron-beam welded together to form the desired multi-attachments of the gas generator pressure chamber. Some prior art inflator housings are shown in U.S. Pat. Nos. 4,711,466; and 4,561,675. However, it has been found that assembling inflator housings by welding is undesirable because it increases their cost of manufacture and weight. Moreover, in some instances it may be advantageous to be able to disassemble a completed housing, for example for quality control, or to replace defective or old chemical agents. However, a welded housing is difficult to take apart."}
{"text": "It has been known for many years that hydrostatic pressures can induce regression of cleavage furrows in dividing eggs (Marsland, D. A. and Brown, D. E. S. J. Cell and Comp. Physiol. 8: 167, 1936). Preliminary investigations with hydrostatic pressures involved a study of the movement of chromosomes by slowing or inhibiting them with hydrostatic pressures, while later studies used pressure to inhibit cell divisions to produce triploid animals, although purely for experimental purposes.\nBeginning with work involving plaice (Purdom, C. E., Heredity 29: 11-24, 1972), many studies have examined methods to induce triploidy in economically important animals. A limited number of these studies have involved the use of hydrostatic pressure to induce triploidy in a variety of fish species while others utilized antibiotics to induce triploidy in shell fish.\nIt is possible to obtain triploids because of the nature of meiosis in lower animals. During meiosis, gametogenic cells undergo a duplication and subsequent division of the chromosomes. In some shellfish, two polar bodies are extruded following fertilization of the eggs. This extrusion can be inhibited by chemical or physical means. Consequently, an extra chromosome set from the polar body may be withheld within the oocyte. This extra set is retained within the egg in conjunction with the set obtained from the sperm, thus producing a triploid embryo.\nTriploids can have practical benefits in aquaculture. These benefits arise from the reproductive sterility that a triploid animal is expected to display. By circumventing the otherwise normal energy expenditure associated with maturation, ripening and spawning of gametes, triploids have been shown to surpass their diploid counterparts in growth and survival (Stanely, J. et al., Aquaculture 37: 147-155, 1984).\nIn the case of Pacific oysters, in particular Crassostrea gigas, limiting the reproductive potential can be especially significant. In general, quality oysters are characterized by high levels of glycogen, which add texture and flavor to the meat. During the breeding season, the market quality of the oyster deteriorates because glycogen is diverted into gamete production instead of being stored. Triploid, and therefore neutered oysters will have reduced potential for maturing.\nConsequently, there exists a need in the art for a method of inducing polyploidy in oysters on a consistent and costs effective basis. The present invention fullfills this need, and further provides other related advantages."}
{"text": "Retailers and merchants involved in electronic commerce continually strive to make it easier for customers to research products and receive information that may assist in purchasing decisions. One way to improve a customer's experience includes providing discussion forums (or discussion boards) that allow customers to exchange information with other people via dialogue in a discussion. The ever-enlarging product catalog makes discussion forums helpful when a customer tries to distinguish between products, seek out impartial advice, learn more about a product or category, share the customer's experience with or impressions of a product, or obtain other information. A merchant with a large catalog of products may offer a large number of different discussion forums on a variety of topics, which may range from very specific topics to general ones. Some forums may be more active and have a greater frequency of responses than other forums. While these forums are useful for customers to share information regarding an item or to make purchasing decisions within a category of items, activity in such forums is not typically analyzed by a merchant or retailer to determine items or categories that are of interest to customers."}
{"text": "1. Field of the Invention\nThe present invention relates to a mounting structure of a canister installed in an automobile.\n2. Description of the Related Art\nIn some automobiles, a spare tire is housed in a rear underfloor region of the vehicle body by use of a chain for hanging the spare tire, and a fuel tank and a canister for adsorbing vaporized fuel are disposed under the floor.\nJapanese Utility Model Application Publication No. S63-105518 discloses a canister mounting structure in which a spare tire is disposed in a substantially central area in the vehicle transverse direction, a fuel tank is disposed in front of the spare tire, an exhaust pipe and a muffler are disposed on one side in the vehicle transverse direction of the spare tire, and a canister is disposed on the opposite side of the spare tire from the muffler."}
{"text": "The present invention relates to a silver halide photographic light-sensitive material, and a method for forming an ultra-high-contrast image, using the same, and particularly to an ultra-high-contrast silver halide photographic light-sensitive material used in the field of graphic arts, and a photographic image-forming method using the same.\nIn order to make reproduction of a line original or a continuous tone original by halftone dots, better in the field of graphic arts, an image-forming system exhibiting a photographic property of ultra-high-contrast (particularly, xcex3=10 or more) is necessary.\nAs a method for obtaining a high-contrast photographic property, a lithographic development manner using the so-called xe2x80x9cinfectious development effectxe2x80x9d has been hitherto used. This manner has the drawback that the developing solution used for it is unstable and thus is difficult to use. An image-forming system, in which development is performed with a developing solution having good storage stability, to obtain an ultra-high-contrast photographic property, has been desired. Examples thereof are disclosed in U.S. Pat. Nos. 4,166,742, 4,168,977, 4,221,857, 4,224,401, 4,243,739, 4,269,922, 4,272,606, 4,311,781, 4,332,878, 4,618,574, 4,634,661, 4,681,836, 5,650,746, and the like. These relate to a system of developing a surface-latent-image-type silver halide photographic light-sensitive material, to which a hydrazine derivative is added, with a stable MQ (Metol-and-Quinol-containing) or PQ (1-phenyl-pyrazolidone-and-hydroquinone-containing) developing solution having a pH of 11.0 to 12.3, to obtain an ultra-high-contrast negative image whose xcex3 is over 10. According to this system, photographic properties of ultra-high-contrast and high sensitivity can be obtained, and a sulfite salt having a high concentration can be added to the developing solution. Therefore, stability of the developing solution against air oxidization is highly improved compared with any conventional lithographic developing solution.\nIn the field of photoengraving, a photographic light-sensitive material having good reproducibility of an original, a stable processing solution, easy replenishment, and the like, are desired, to cope with variation and complication of printed matters.\nAn original, particularly in the step of photographing an original comprising line images, is made by combining typeset characters, handwritten characters, illustrations, a photograph converted to halftone dots, and the like. In one original, therefore, image portions having different densities and line widths coexist being intermixed. A process camera, a photographic light-sensitive material, or an image-forming method for finishing such an original with high reproducibility, are intensely desired. In the plate-making of a catalogue or a large-sized poster, enlarging or scaling down of a dot photograph is widely carried out. In plate-making in which dots are enlarged, lines are made large, to give images in which dots fade. In plate-making in which dots are scaled down, the number of lines per inch becomes larger than that of an original, to give images in which dots are very fine. Accordingly, in order to keep halftone dot reproducibility good, an image-forming method that can give wider latitude is desired.\nTo improve reproduction of the original, JP-A-3-39952 (the term xe2x80x9cJP-A-xe2x80x9d as used herein means an unexamined published Japanese patent application), JP-A-3-174143, or the like disclose an ultra-high contrast silver halide light-sensitive material having a multilayer structure, which includes a layer containing a redox compound that releases a development inhibitor by oxidation, and a light-sensitive silver halide emulsion layer containing a hydrazine derivative. In this method, reproduction of an original is remarkably improved.\nHowever, the improving of original-reproducibility is an eternal theme for graphic art light-sensitive materials. Thus, a system giving a better original-reproducibility is intensely desired.\nThe above method, in which a surface-latent-image-type silver halide photographic light-sensitive material, to which a hydrazine derivative is added, is processed in a developer having a pH of 11.0 to 12.3, makes it possible to improve the stability of a developer by using a high-density sulfurous acid preservative. However, in order to obtain an ultra-high contrast photographic image, it is necessary to use a developer having a relatively high pH value. In this case, the developer tends to be oxidized by air, and it is therefore necessary to replenish a large amount of a developer. In view of this, various ideas, in which an ultra-high contrast photographic image-forming system utilizing nuclei formation developing of a hydrazine compound is realized using a developer having a lower pH, have been attempted.\nIn order to obtain an ultra-high contrast image by using a stable developer having a pH of lower than 11.0, a method using a highly active hydrazine-type nucleating agent, and a nucleation accelerator, is disclosed in U.S. Pat. No. 4,269,929 (JP-A-61-267759), U.S. Pat. No. 4,737,452 (JP-A-60-179734), U.S. Pat. Nos. 5,104,769, and 4,798,780; JP-A-1-179939 and JP-A-1-179940; U.S. Pat. Nos. 4,998,604 and 4,994,365, and JP-A-8-272023.\nHowever, the above method cannot be said to be satisfactory in the point of original-reproducibility in the case of using the aforementioned redox compound that releases a development inhibitor by oxidation. There is a strong need to develop light-sensitive materials having higher original-reproducibility. Also, in view of environmental problems, which have been greatly emphasized in recent years, there is a strong need to develop light-sensitive materials that cause less fluctuation in the photographic properties, and that also cause less fluctuation in original-reproducibility, with change in the composition of a developer, to decrease developer waste.\nOn the other hand, Japanese Patent No. 2,676,439 discloses a hydrazine-type redox compound containing a sulfo group, and JP-A-9-269553 discloses a hydrazine-type redox compound having in its molecule two sulfonamido groups, and a redox compound having phenol as its moiety. However, these compounds make it impossible to give a silver halide photographic light-sensitive material having good original-reproducibility, and suppressed black-spots (black peppers), high sensitivity, and high contrast, by use of a developing solution having a low pH.\nThe present invention is a silver halide photographic light-sensitive material, which comprises at least one compound represented by the following formula (1): \nwherein PUG represents a nitrogen-containing heterocyclic group which is bonded to the adjacent carbonyl group through a nitrogen atom, the heterocyclic group is a heterocyclic group which directly or indirectly has at least nitro group as its substituent(s) and can give a pKa value of 9 to 11, ph represents a phenyl group and directly or indirectly has at least one dissociating group as its substituent(s), Z represents a substituent which is capable of substitution on the benzene ring, and p is an integer of 0 to 4.\nFurther, the present invention is a method of processing a silver halide photographic light-sensitive material, which comprises the step of developing the silver halide photographic light-sensitive material described above with a developing solution having a pH of 9.0-11.0 after imagewise exposure to light.\nOther and further features and advantages of the invention will appear more fully from the following description."}
{"text": "The present invention relates to a picture frame for framing pictures, posters, prints, signs and the like. More particularly, it relates to such a picture frame for the inexpensive, yet aesthetically pleasing and practical framing of a wide variety of pictures for both residential and commercial use.\nIn my earlier patent for a picture frame, U.S. Pat. No. 4,669,209, I have described a picture frame which is very inexpensive and so easy to assemble that anyone can do so with aesthetic and precise results. The picture frame is designed to eliminate the need for a picture glass or other transparent covering and semi-rigid backing, normally used for most picture frames. It was specifically intended for framing posters, which tend to be relatively large, and prints. Artistic decor utilizing posters and large prints has come into vogue in recent years. Where such poster-type pictures are inexpensive, it is often times not worthwhile to provide the same with a glass covering and a semi-rigid backing, together with a frame, for the entire assembly since these various extra elements often cost far in excess of the original poster to be framed. Thus, on many occasions, one sees attractive poster-like pictures which are tacked or otherwise fastened to a wall for the purpose of exhibiting the picture and maintaining it flat. Such securement of a poster-like picture to a wall not only causes damage to the wall but very often tends to be unsightly, in that tacks or tape, or other objects are utilized in fastening the picture to the wall.\nTo overcome this problem, my earlier patent provided a picture frame for poster-like pictures which included at least a top picture edge support, and preferably a bottom picture edge support, and fastening elements which fasten the respective edges of the poster-like picture frame to the top and bottom supports. Each picture edge support includes a flat back plate which extends substantially the length of the edge of the picture being framed and a substantially C-shaped portion or member, one leg of which is connected to the outer edge of the back plate, while the other leg extends downwardly toward the flat back plate to define a gap between the terminus thereof and the back plate. The picture edge support is preferably formed of extruded plastic. The fastening elements are also generally C-shaped and resiliently engaged around the generally C-shaped member of the edge support with the inside leg extending into the gap beyond the terminus of the inside leg of the C-shaped member and the flat back plate, so as to press the edge of the poster-like picture inserted into the gap against the back plate of the edge support. The invention also contemplates the provision of edge supports for each side edge of the poster-like picture, wherein the fastening elements are in the form of elbow-shaped corner pieces, which engage the adjacent ends of adjacent edge supports so as to secure or fasten the poster-like picture to the flat back plates of the supports at the corners of the picture. While quite satisfactory in use, I have found that for certain applications a more rigid or stiffer picture frame is needed. In addition, the picture frame is not suitable for joining multiple pictures or posters together. It is also not particularly suitable for double-sided signs, such as may be often used in store windows and the like, where it is desired that the customer can read the same either on the outside of the store or on the inside thereof."}
{"text": "A common method for forming parts from plastic involves the extrusion of material from an extruder. The extruded plastic, or extrudate, however, remains in the molten state immediately upon exiting the extruder. As a result, the extrudate is subject to undesirable deformation until it cools sufficiently to become solid.\nOver the years, efforts have been made in the plastics industry to control the size and shape of extruded plastic profiles and extruded plastic tubing as they cool and harden. One widely used method of cooling extruded profiles consists of passing the profile onto an air rack after it exits the extrusion die. Compressed air is directed or blown onto the extrudate to facilitate cooling while the air rack maintains the size and shape of the profile as it solidifies. Another method of maintaining the size and shape of the extruded plastic involves the spraying or cascading of water onto the extrudate after it exits the extruder, or immersing the extrudate in a water bath as it exits the extruder. Spraying and water tanks of this type have been commonly used in the cooling of extruded plastic profiles. However, tanks of this type are often required to be of substantial length to insure that the extrudate is sufficiently cooled and hardened before exiting the tank. As a result, commercially available spraying and water tanks are typically of substantial length, twelve feet or longer, to ensure sufficient cooling.\nTypically, alternative methods include passing the extrudate through a sizing and cooling tool or die, which is immersed and contained within a water tank, to help maintain the size and shape of the extrudate as it cools. Effective sizing tools or dies are often quite complex and require precision machining. Often, a sizing tool or die is comprised of a combination of precision machined tapered components which impart the desired size and shape to the extrudate. Various types of sizing or calibrating tools appear in Kessler, U.S. Pat. No. 4,181,487 and Gauchel, U.S. Pat. No. 4,411,613. Due to the precision machining required, effective sizing tools or dies are often quite expensive. Furthermore, sizing tools or dies are often of substantial length, typically six to eight inches or longer.\nIt is has been found that water in the spraying or water tanks exerts a pressure on the extrudate which can cause unwanted deformation of the extrudate as it cools. Extruded tubing and other hollow extruded articles are particularly susceptible to deformation during the cooling process. To help eliminate the deformation caused by the water in the cooling tanks, it has been found advantageous to place the water under a negative pressure. This has been accomplished through the use of a vacuum pump to create a less than atmospheric pressure in the water of the cooling tank. The vacuum within the cooling tank coupled with the higher pressure within the hollow portion of the extrudate helps to maintain the shape of the extrudate during the cooling process. Sizing tools have also been adapted with vacuum lines to enable vacuum pressure to be applied to the extrudate as it passes through the sizing tool.\nVacuum sizing tanks with complex vacuum pressure control systems have been developed, such as ones disclosed in DeCoursey, U.S. Pat. No. 5,008,051 and Preiato, U.S. Pat. No. 5,340,295, the disclosure of which are herein incorporated by reference. The use of elaborate vacuum sizing tanks with carefully controlled vacuum pressures are particularly useful where precision tubing is involved. However, commercially available vacuum sizing tanks of this type are expensive and of substantial length to accommodate a variety of applications. Furthermore, carefully controlled vacuum pressure is often unnecessary for profile parts and low precision hollow extruded parts.\nThe above-described water tanks typically consist of an enclosed trough or close-walled tank holding water the entire length of the tank, rendering further modification or manipulation of the extrudate impractical, after it passes through the sizing tool or exits the tank. Vacuum sizing tanks are structurally enclosed to hold the vacuum pressure, which also renders further manipulation or modification of the extrudate after passing through the sizing tool impractical.\nTherefore, it is an object to the present invention to provide a versatile vacuum sizing system wherein the vacuum sizing tank is substantially reduced in length, and wherein said apparatus is useful for a variety of applications. It is a further object of the invention to reduce the high cost of lengthy precision machined sizing tools by substantially decreasing the length of the sizing tool or die. It is a further object of the invention to provide a vacuum sizing system wherein the shape of the extrudate can be further manipulated or modified in an ambient air environment after passage through a vacuum sizing tank."}
{"text": "1. Field of the Invention\nThis invention relates to kits for use with oil level indicators of fuel tanks, to provide at a remote point, indications of fuel level in the tank.\n2. Description of the Related Art Including Information Disclosed Under 37 CFR .sctn..sctn.1.97-1.99\nFuel tanks located in the basement areas of dwellings are usually provided with a simple mechanical level indicator that is visible at the top of the tank. It is often inconvenient or undesirable for the dweller of the building to descend to the basement area to continually check the oil level to insure against the tank running out of oil."}
{"text": "1. Field of the Invention\nThe present invention relates to devices and methods for clamping the leads of a lead frame while the leads are being wire-bonded to bond pads on an integrated circuit die.\n2. State of the Art\nIntegrated circuit (IC) dice or “chips” are small, generally rectangular IC devices cut from a semiconductor wafer, such as a silicon wafer, on which multiple ICs have been fabricated. IC dice are typically packaged to protect them from corrosion by attaching them to lead frames using a solder or epoxy, electrically connecting them to leads in the lead frames, and then encapsulating the resulting die and lead frame assemblies in plastic die packages. The leads in the encapsulated lead frames protrude from the die packages and terminate in pins so the dice electrically connected to the leads can communicate with circuitry external to the die packages through the pins.\nIC dice are typically electrically connected to the leads in their associated lead frames in a process known as “wire-bonding.” During wire-bonding of an IC die to a lead frame, a wire-bonding tool, such as an ultrasonic, thermosonic, or thermo-compression tool, bonds lengths of fine wire to bond pads on the die and to bonding sites on the leads in the lead frame in order to connect the bond pads and the leads. This process is well-known, and is described in detail in U.S. Pat. Nos. 4,030,657, 4,361,261, 4,527,730, 4,600,138, 4,603,803, 4,765,531, 4,778,097, 5,148,959, 5,217,154, 5,251,805, 5,307,978, 5,421,503, 5,425,491, 5,445,306, 5,465,899, 5,486,733, 5,513,792, and 5,516,023.\nIn order to ensure proper wire-bonding, lead frame leads are typically clamped during wire-bonding to stabilize them so a solid bond can be formed between the leads and the fine wire used to wire-bond them to an IC die. Without clamping, the leads can move during wire-bonding and cause unreliable bonds to be formed between the leads and the fine wire.\nConventionally, all the leads in a lead frame are clamped at once in a single plane during wire-bonding, as shown in U.S. Pat. Nos. 3,685,137, 5,035,034, and 5,193,733. By clamping all the leads of a lead frame at once, the conventional clamping method holds to a relative minimum the time it takes to wire-bond an IC die to the lead frame.\nWhile holding wire-bonding time to a minimum is advantageous, the conventional clamping method can also be problematic with certain lead frames. For example, some lead frames, such as the Leads-Over-Chip lead frame shown in U.S. Pat. No. 4,862,245, have leads which extend closer to the bond pads of IC dice attached to the lead frames than other leads such that fine wires bonded between the IC dice and these other leads must arch over the leads which are closer to the bond pads on the dice. Since fine wires must arch over the leads closer to the bond pads, the closer leads cannot be clamped at the same time as all the other leads are clamped. Instead, the closer leads conventionally remain unclamped, or only partially clamped, during wire-bonding, resulting in potentially unreliable bonds to those leads.\nAlso, because the conventional clamping method clamps all the leads of a lead frame at once in a single plane, it cannot adequately clamp leads of a lead frame which must be clamped in different planes. For example, as shown in U.S. Pat. No. 4,943,843, some Tape-Under-Frame lead frames have leads which extend underneath an IC die and then project beyond the footprint of the die so they may be wire-bonded to bond pads on the die. These leads are typically left unclamped during wire-bonding, because they are not clampable in accordance with the conventional clamping method in the same plane as the remainder of the leads. As a result, these leads have potentially unreliable wire bonds. In a similar fashion, the conventional clamping method cannot adequately clamp multi-layer lead frames, such as that shown in U.S. Pat. No. 5,291,061.\nWire-bonding clamps, such as that shown in U.S. Pat. No. 4,821,945, have been developed to alleviate some of the problems of the conventional clamping method by clamping leads individually. However, these clamps can dramatically increase the time it takes to wire-bond an IC die, because the clamp must be sequentially indexed to each individual lead.\nTherefore, there is a need in the art for a device and method for clamping lead frame leads during wire-bonding that can clamp leads in different planes without significantly increasing the time it takes to wire-bond an IC die."}
{"text": "Elastomeric meltblown webs have been proposed for use in a variety of products including composite fabrics including hydroentangled fabrics; in diapers, training pants and other personal hygiene products in which stretch and conformability to body shapes are considered important. Fully hydrogenated (saturated) diblock and/or triblock copolymers and mixtures thereof based on polystyrene blocks and poly(ethylene-butylene) blocks have been the subject of considerable attention for producing meltblown elastomeric webs because of their high temperature stability and their ability to produce meltblown webs with desirable properties.\nCommercially available polystyrene-(ethylene-butylene) diblock and triblock copolymers include the KRATON-G resins commercially available from Shell Chemical Company. Because of the high viscosities associated with these resins, the manufacturer's literature suggests blending of the resins with certain relatively low molecular weight materials. The blending of such materials with the KRATON resins can reduce the processing temperatures, thereby minimizing the degradation of the materials, or can reduce melt processing viscosities, thereby enabling throughputs to be increased at lowered pressures in extrusion processes, such as meltblowing processes. The Shell literature teaches that the lower molecular weight materials which are useful in blends include those which are compatible with the polystyrene (PS) segments of the copolymer, and materials which are compatible with the ethylene-butylene (EB) segments. Materials which are compatible with the (PS) segments include polystyrene and poly(methylacrylate) while polyolefins are compatible with the (EB) segments.\nU.S. Pat. No. 4,663,220 to Wisneski and U.S. Pat. No. 4,692,371 to Morman disclose the preparation of meltblown webs from blends of saturated (PS)-(EB) diblock and triblock elastomers together with polyolefin resins. However, the preparation of meltblown webs at high throughput rates using these blends can result in processing difficulties rendering the high throughput meltblowing process uneconomical.\nU.S. Pat. No. 4,323,534 to Des Marais discloses the use of fatty acids or fatty alcohols as plasticizers useful in the meltblowing of KRATON G, fully saturated elastomers. More recently, U.S. Pat. No. 4,892,203 to Himes discloses blends of the fully saturated KRATON G-type resins plasticized with anionically polymerized styrene or alpha-methyl styrene or their copolymers, or hydrogenated polystyrene. Optionally, a microcrystalline wax may also be added.\nU.S. Pat. No. 4,874,447 to Hazelton discloses a method for preparing a nonwoven web from a blend comprising (i) an elastomeric copolymer of an isoolefin and a conjugated diolefin, and (ii) a thermoplastic olefin polymer resin. The elastomers (i) disclosed include copolymers of styrene and butadiene, but none of the fully hydrogenated block copolymers of the KRATON G-type are disclosed. A wide range of thermoplastic resins are disclosed as component (ii), including polyolefins, such as polyethylene, polypropylene, polybutylene, polypentene, copolymers of ethylene and propylene, copolymers of ethylene with unsaturated esters of lower carboxylic acids including copolymers of ethylene with vinylacetate or alkyl acrylates, and the like. However, the unsaturated block copolymers lack the high temperature stability of the saturated block copolymers, and thus elastomeric webs from these materials or blends of these materials can be more difficult to process.\nU.S. Pat. No. 4,769,279 to Graham discloses meltblown webs formed from blends of ethylene-acrylic copolymer or ethylene-vinylacetate blended with a second fiber-forming polymer such as a polyolefin. However, the elastomeric webs formed from blends based on ethylene-acrylic copolymers and/or ethylene vinylacetate copolymers, as the elastomeric material, have only limited stretch and recovery properties.\nDespite substantial effort and experimentation in the art, only a limited number of elastomeric materials have been used with any substantial commercial success to produce elastomeric webs. Moreover, various processing difficulties are still encountered when attempts are made to produce meltblown elastomeric webs at relatively high throughput rates."}
{"text": "The present invention relates generally to hard-bodied artificial fishing lures and specifically to lures with retractable hooks which make them xe2x80x9cweedlessxe2x80x9d or xe2x80x9csnagless.xe2x80x9d\nArtificial fishing lures are generally designed to mimic the look and action of some species of bait fish. Manufacturers have been building lures of various types for many years. The bodies have been manufactured from various substances including balsa and other woods, red cedar, plastic, and metal. The shapes are generally designed to imitate bait fish of some bait species. The lures come in almost infinite colors and surface finishes.\nAlmost since the invention of the first lure of this type (xe2x80x9ccrank baitxe2x80x9d, xe2x80x9ctop waterxe2x80x9d lures, or xe2x80x9clipless crank batexe2x80x9d) manufacturers have used hooks with multiple barbs to xe2x80x9chookxe2x80x9d a fish when a strike occurs. Most lures use treble hooks with three distinct barbs, and are suspending from belly or bottom side of the lures.\nThe dangling hooks on most lures, while effective in securing a catch once a strike occurs, are problematic in other contexts. The dangling hooks are easily entangled in brush, rocks, moss and surface/sub-surface debris, often resulting in loss of an expensive lure.\nManufacturers have tried to make their lures more xe2x80x9cweedlessxe2x80x9d by adding various hook guarding mechanisms. However, hook guards are difficult to install on lures employing multiple, dangling treble hooks. In recent years some manufacturers have attempted to make their lures more weedless by insuring a xe2x80x9cnose downxe2x80x9d attitude for the lures as they are drawn through the water, and this approach does provide some amount of protection from snags. In top-water applications, on the other hand, some work has been done using soft bodies which protect the hook until the body is compressed in the fishes mouth.\nOther approaches to achieving xe2x80x9cweedlessxe2x80x9d luresxe2x80x94those with hidden, or retracted hooksxe2x80x94have their own problems or shortcomings. Most, if not all, such designs, cannot be produced as a single piece casting or two piece casting using current molding and machining technology. Many require intricately designed and machined parts to function, and involved complicated assemblies with many moving pieces.\nDesigns with retracted, strike-deployed hoods include that of Zeman (U.S. Pat. No. 3,646,699) which employs a single hook that is activated to the hooking position by use of a spring. Braden, in U.S. Pat. No. 4,765,084 and in U.S. Pat. No. 5,878,524, teaches a mechanism with more than fifty distinct parts. Assembly of those mechanisms into an assembled body would prove costly, tedious, and the end product unreliable. Durability is a problem associated with, or arising direction from some of the just-described designs. The known designs employ springs, cams, levers, and other actuating devices that are soon destroyed by rust, corrosion and breakage.\nAnother problem with retracted hook designs for xe2x80x9cweedlessxe2x80x9d lures is that of pre-mature hook deployment. Previous designs employ complicated and delicate mechanisms to hold the hook in the retracted position. In U.S. Pat. No. 3,665,634 Baud employs a cam that sets into a detent in the hook body. The force of a cast will dislodge the cam and cause early deployment in normal use. Hameen and Antiua in U.S. Pat. No. 3,816,953 employ a complicated spring and lever mechanism. In somewhat contradictory fashion, another problem is that of keeping a hook deployed, once extended. In U.S. Pat. No. 4,528,770, McDiarmand employs a complicated spring mechanism to deploy the hook, but a spring strong enough to hold the hook in place may cause the hook to fail to deploy upon a strike.\nYet another limitation of existing, xe2x80x9cweedlessxe2x80x9d designs relates to the shape of the end product. To provide room for complicated moving parts the lure bodies are often fattened, elongated, and bulged and the end product has no real resemblance to bait.\nFinally, there is the problem with retracted hook lures of depending on non-standard parts. In U.S. Pat. No. 4,528,770, McDiarmand employs a hook that is not standard industry design, as does Bates in U.S. Pat. No. 4,873,781 Bates.\nIn summary, the existing approaches to making lures xe2x80x9cweedlessxe2x80x9d are only somewhat effective, and true xe2x80x9cweedlessxe2x80x9d construction has not yet been attained.\nLures equipped with treble hooks, whether intended to be xe2x80x9cweedlessxe2x80x9d or not, present several additional drawbacks and problems. A first problem has to do with safety. Once a fish takes the lure into its mouth and becomes ensnared in the hooks, a fisherman is faced with the extraction of the hooks once the fish is landed or netted. To extract the hooks a fisherman must restrain the fish while pulling the hook(s) from the fishes mouth. Because most lures contain multiple treble hooks, any thrashing or movement of the fish during the extraction can result in the xe2x80x9cfreexe2x80x9d hooks (hooks not caught in the fishes mouth) becoming embedded in the fisherman. Another problem relates to fish being often able to dislodge the treble hooks. This is particularly true of species which are known to leap and thrash their heads back and forth. This is less of a problem when single barb hooks are used, because of their ability to provide a solid hook-up and retention of the fish due to the larger size of the hook and barb.\nAn advantage of a lure with a single hidden hook relates to storage. Most hard bodied baits have to be stored in separate compartments to prevent tangling. When combined with other treble-hooked lures, hooks become entangles and time is wasted in the untangling process. If the lures are stored in separate compartments, additional storage space and equipment is needed. Lures with a hidden hook could be stored with multiple lures to a compartment with no difficulties.\nIn view of the foregoing, it is an object of the present invention to provide an improved, artificial fishing lure.\nIt is another object of the present invention to provide an improved artificial fishing lure which more effectively than existing design, resists entanglement with weeds and the like during use.\nIt is another object of the present invention to provide an improved artificial fishing lure of retracted hook design which is more reliable than existing entanglement-resistant (xe2x80x9cweedlessxe2x80x9d) lure design.\nIt is another object of the present invention to provide an improved artificial fishing lure of retracted hook design which is less mechanically complex than existing entanglement-resistant (xe2x80x9cweedlessxe2x80x9d) lure design.\nIt is another object of the present invention to provide an improved artificial fishing lure of retracted hook design which is more escape-resistant than existing entanglement-resistant (xe2x80x9cweedlessxe2x80x9d) lure design.\nIt is another object of the present invention to provide an improved artificial fishing lure of retracted hook design which is less hazardous to users than existing entanglement-resistant (xe2x80x9cweedlessxe2x80x9d) lure design.\nIn satisfaction of these and related objects, the present invention provides an improved hard bodied lure with a single retractable hook, adaptable to multiple configurations, which can be readily manufactured using current molding and machining technology and standard industry parts, and which provides true weedless/snagless operation/use Since most game fish have an affinity for underwater structures such as brush and rocks, being able to fish the lure around and through this type of structure is extremely advantageous.\nAs will be clear from a review of the Detailed Description, the present design is simple, effective, and addresses each of the problems enumerated above in the Background. The present design guards the hook while the lure is being fished, and provides for deployment only when the fish strikes. This is accomplished via a cavity in the body of the lure. The cavity is large enough to house the hooking point of the single attached hook. The hook retention mechanism keeps the hook from deploying until the fisherman desires the hook to deploy and creates a source of biasing force to deploy the hook upon demand.\nThe present design addresses safety concerns described above by using a single barb hook, and in like manner addresses the dislodging problem of treble hooks. The simplicity of the mechanism allows for outer contours which meet whatever shapes are desirable for optimum fish attraction. The bodies of lures of the present design may be readily produced with molding technology employed by the industry today as a single piece or two piece casting, or as a single or two piece machined unit of wood, metal, or plastic using current machining and lathe technologies such as CNC machines. The present design is, mechanically speaking, very simple, and not prone to mechanical failure nor to rust or other deterioration. This same simplicity is beneficial from a cost of manufacturing standpoint. Other benefits of embodiments of the present invention will be apparent from an understanding of the specification and associated drawings."}
{"text": "I. Field of the Invention\nThis invention relates to a process for producing fibres composed of, or coated with, carbides or nitrides of such elements as e.g. silicon and boron.\nII. Description of the Prior Art\nContinuous and short fibres made of SiC and similar carbides and nitrides are very useful for reinforcing a variety of materials such as metals and ceramics in order to improve the toughness or high temperature performance of such materials. The usefulness of the fibres arises not only from their inherent hardness and tensile strength, but also from the fact that the fibres generally are stable in oxidizing atmospheres and at high temperatures and tend not to react with the ceramic or metal matrix.\nFrom these points of view, fibres of this kind are superior to carbon fibres often used for the same purposes. For example, when carbon fibres are used to reinforce aluminum metal, enormous pressures must be used to cause the molten aluminum to impregnate the fibres because carbon is not wetted by the molten aluminum. Moreover, the carbon reacts with the aluminum to form Al.sub.4 C.sub.3 at the interface and the resulting composite is of low strength. Attempts have been made in the past to overcome these problems by coating carbon fibres with materials such as metals (e.g. Ni) or carbides or borides, etc. in order to make the fibres wettable and to make their surfaces inert. The coating methods attempted have included plasma spraying, chemical vapour deposition, sputtering, etc. but there has been limited success due to poor penetration of the coating materials into fibre bundles, high costs and unacceptable losses of coating material etc.\nDespite the superiority of fibres made of carbides and the like as mentioned above, the conventional fibres of this kind have a number of problems For example, SiC fibres are often manufactured by degrading a polycarbosilane type of precursor, but the high temperature stability of such fibres is low because they undergo a crystallization change at temperatures of about 1300.degree.-1400.degree. C. and thus tend to fall apart at these temperatures. It is also difficult to manufacture long or continuous fibres of carbides and the like, and although it has been reported that continuous SiC fibres can be made by heating a precursor fibre containing SiO.sub.2 or Si particles distributed in a carbon precursor such as pitch, attempts to duplicate such a process have not been successful. SiC particles are indeed formed, but they do not sinter together even at temperatures exceeding 2000.degree. C. and thus they do not retain the original fibre shape.\nIt has been known that suboxides such as SiO can be formed by reacting SiO.sub.2 and C or SiO.sub.2 and Si etc. and reacted with a graphite surface to form a coating. For example JP 58104078-A, JP 540090216-A etc. refer to the reaction of SiO(g) with graphite to form a coating on the surface. The SiO(g) is produced in a separate reactor and transported to the graphite surface. This is feasible at temperatures above 1800.degree. C. as at lower temperatures SiO can re-convert to SiO.sub.2 and Si (or, in presence of CO, to SiO.sub.2 and C). Coating of a carbon fibre mass by this method results in two problems:\na) due to improper infiltration, the coating formed is non uniform; and\nb) as the reaction is conducted at high temperatures, the coatings are very granular and highly porous.\nAccordingly, there is a need for an improved process for producing fibres made of, or coated with, carbides or nitrides, which process overcomes some or all of the disadvantages mentioned above."}
{"text": "1. Field of the Invention\nEmbodiments of this invention are directed to the field of intravascular medical devices, and more particularly to the field of catheters such as angioplasty, neurological and guide catheters, among others, which may be used in various medical procedures such as percutaneous transluminal angioplasty (PTA), percutaneous transluminal coronary angioplasty (PTCA) as well as in procedures involving the placement of medicines and medical devices within the body.\nSome embodiments of the invention are directed to all forms of catheters which may be advanced through a body lumen or vessel. Some examples of catheters are over-the-wire (OTW) catheters, such as are described in U.S. Pat. No. 5,047,045; single-operator-exchange (SOE) balloon catheters, such as are described in U.S. Pat. Nos. 5,156,594 and 5,549,552. Other examples of catheters which may incorporate the unique features of the present invention may include rapid-exchange catheters, guide catheters, etc.\n2. Description of the Related Art\nIntravascular diseases are commonly treated by relatively non-invasive techniques such as PTA and PTCA. These angioplasty techniques typically involve the use of a balloon catheter. In these procedures, a balloon catheter is advanced through the vasculature of a patient such that the balloon is positioned proximate a restriction in a diseased vessel. The balloon is then inflated and the restriction in the vessel is opened. In other uses a catheter may be used to delivery an endoprosthesis such as a stent, graft, stent-graft, vena cava filter or other implantable device or devices herein after collectively referred to as a stent or stents. Where a stent is to be delivered into a body lumen the catheter may include one or more inflatable portions or balloons. Typically, the stent is retained in the predelivery state about the catheter shaft, or a portion thereof such as a balloon, by crimping and/or through the use of a retaining mechanism such as sleeve, sheath or sock.\nBalloons and balloon catheters may be particularly useful for the delivery of expandable, implantable medical devices such as stents, grafts, stent-grafts, vena cava filters, hereinafter referred to cumulatively as stents. Stents and catheters used in their delivery are commonly used and as such their structure and function are well known.\nA stent is a generally cylindrical prosthesis introduced via a catheter into a lumen of a body vessel in a configuration having a generally reduced diameter and then expanded. In its expanded configuration, the stent supports and reinforces the vessel walls while maintaining the vessel in an open, unobstructed condition.\nIn order to properly position a stent and/or balloon within a body lumen, the catheter must be advanced through the narrow confines of the body. Typically the balloon and/or stent is located near the distal end of the catheter. In order to advance the distal end of most prior catheters further in to a body lumen, the inner shaft or catheter member is utilized to transmit force to the distal end. However, the inner is typically soft and flexible which often results in poor transmission of push.\nThe present invention, in accordance with the various embodiments presented herein, addresses the shortcoming of poor push transmission common to many catheters.\nWithout limiting the scope of the invention, a brief summary of various embodiments of the invention is set forth below. Additional details of the summarized embodiments of the invention and/or additional embodiments of the invention may be found in the Detailed Description of the Invention below.\nA brief abstract of the technical disclosure in the specification is provided as well for the purposes of complying with 37 C.F.R. 1.72.\nThe entire content of all of patents or other references listed within the present patent application are incorporated herein by reference."}
{"text": "Large quantities of methane, the main component of natural gas, are available in many areas of the world, and natural gas is predicted to outlast oil reserves by a significant margin. However, most natural gas is situated in areas that are geographically remote from population and industrial centers. The costs of compression, transportation, and storage make its use economically unattractive.\nTo improve the economics of natural gas use, much research has focused on methane as a starting material for the production of higher hydrocarbons and hydrocarbon liquids. The conversion of methane to hydrocarbons is typically carried out in two steps. In the first step, methane is reformed with water to produce carbon monoxide and hydrogen (i.e., synthesis gas or syngas). In a second step, the syngas is converted to hydrocarbons, for example, using the Fischer-Tropsch process to provide fuels that boil in the middle distillate range, such as kerosene and diesel fuel, and hydrocarbon waxes. Present day industrial use of methane as a chemical feedstock typically proceeds by the initial conversion of methane to carbon monoxide and hydrogen by either steam reforming, which is the most widely used process, or by dry reforming. Steam reforming proceeds according to Equation 1.CH4+H2O CO+3H2  (1)Although steam reforming has been practiced for over five decades, efforts to improve the energy efficiency and reduce the capital investment required for this technology continue.\nThe partial oxidation of hydrocarbons, e.g., natural gas or methane is another process that has been employed to produce syngas. While currently limited as an industrial process, partial oxidation has recently attracted much attention due to significant inherent advantages, such as the fact that significant heat is released during the process, in contrast to the steam reforming processes, which are endothermic. Partial oxidation of methane proceeds exothermically according to the following reaction stoichiometry:CH4+1/202 CO+2H2  (2)\nIn the catalytic partial oxidation processes, natural gas is mixed with air, oxygen or oxygen-enriched air, and is introduced to a catalyst at elevated temperature and pressure. The partial oxidation of methane yields a syngas mixture with a H2:CO ratio of 2:1, as shown in Equation 2. This ratio is more useful than the H2:CO ratio from steam reforming for the downstream conversion of the syngas to chemicals such as methanol and to fuels. Furthermore, oxidation reactions are typically much faster than reforming reactions. This makes possible the use of much smaller reactors for catalytic partial oxidation processes. The syngas in turn may be converted to hydrocarbon products, for example, fuels boiling in the middle distillate range, such as kerosene and diesel fuel, and hydrocarbon waxes by processes such as the Fischer-Tropsch synthesis.\nThe selectivities of catalytic partial oxidation to the desired products, carbon monoxide and hydrogen, are controlled by several factors, but one of the most important of these factors is the choice of catalyst composition. Difficulties have arisen in the prior art in making such a choice economical. Typically, catalyst compositions have included precious metals and/or rare earths. The large volumes of expensive catalysts needed by the existing catalytic partial oxidation processes have placed these processes generally outside the limits of economic justification.\nA number of process regimes have been described in the literature for the production of syngas via catalyzed partial oxidation reactions. The noble metals, which typically serve as the best catalysts for the partial oxidation of methane, are scarce and expensive. The more widely used, less expensive, catalysts have the disadvantage of promoting coke formation on the catalyst during the reaction, which results in loss of catalytic activity. Moreover, in order to obtain acceptable levels of conversion of gaseous hydrocarbon feedstock to CO and H2 it is typically necessary to operate the reactor at a relatively low flow rate, or space velocity, using a large quantity of catalyst. For successful operation at commercial scale, however, the catalytic partial oxidation process must be able to achieve a high conversion of the methane feedstock at high gas hourly space velocities, and the selectivity of the process to the desired products of carbon monoxide and hydrogen must be high. Such high conversion and selectivity must be achieved without detrimental effects to the catalyst, such as the formation of carbon deposits (“coke”) on the catalyst, which severely reduces catalyst performance.\nAs a result, substantial effort has been devoted in the art to the development of economical catalysts allowing commercial performance without coke formation. Not only is the choice of the catalyst's chemical composition important, the physical structure of the catalyst and catalyst support structures must possess mechanical strength, in order to function under operating conditions of high pressure and high flow rate of the reactant and product gasses.\nOf the methods that employ catalysts for oxidative conversion of methane to syngas, typically catalytic metals are dispersed throughout a porous support. The porous support provides longitudinal channels or passageways that permit high space velocities with a minimal pressure drop. In ideal conditions, the catalytic metals are dispersed throughout the porous channels and promote further conversion.\nThere is a continuing need for better, more economical processes and catalysts for the catalytic partial oxidation of hydrocarbons, particularly methane, or methane containing feeds, in which the catalyst retains a high level of activity and selectivity to carbon monoxide and hydrogen under conditions of high gas space velocity and elevated pressure.\nAs is known, in a diffusion-limited heterogeneous process, the contribution of internal active surface areas is defined with an effectiveness factor (η) that has a value ranging from 0 to 1. The effectiveness factor (η) is determined by the intrinsic reaction kinetics and mass transfer coefficients, and can be expressed as a function of the Thiele modulus (Φ), which in turn is given by Equation 3:Φ=L√{square root over (r/Deff)}  (3)where L is the diffusion distance (e.g. radius for sphere); r is the intrinsic reaction rate, and Deff is the effective diffusion coefficient of the reactants. The effectiveness factor (η) is related to the Thiele modulus as shown in FIG. 1.\nAs indicated in FIG. 1, the effectiveness factor (η) can be increased by decreasing the Thiele modulus (Φ). In porous materials, the effective diffusion coefficient (Deff) is a function of gas property and material porosity (pore volume and pore size distribution). This relationship is qualitatively illustrated in FIG. 2.\nAs indicated in FIG. 2, the effective diffusion coefficient (Deff) can be increased by increasing pore size. For example, for pores in the range of 1-100 nm, increasing pore size by 10-fold can roughly increase Deff by an order of 10, which decreases Φ value 3-fold. As the result, the effectiveness factor (η) can be increased substantially."}
{"text": "Digital halftoning is a technique employing digital image processing to produce a halftone output image from a continuous-tone input image. In the digital halftoning technique, a continuous-tone image is sampled, for example, with a scanner or digital camera, and the samples are digitized and stored in a computer. The digitized samples (or “pixels”) include discrete code values having Ni possible input levels (typically Ni=256, with code values ranging from 0 to 255). To reproduce this image with high quality on an output device which can print No output levels, where No<Ni, it is necessary to produce the sensation of intermediate gray levels by suitably distributing the printed dots in the output image using some form of halftoning technique. One example of an output device requiring halftoning is an inkjet printer, which ejects a volume of ink at each output pixel to reproduce the image. The number of output levels (No) of the printer is the number of possible volumes deposited at each pixel (including 0), and is typically in the range of 2-32.\nOne prior art method of digital halftoning is known as error diffusion. FIG. 1 shows a block diagram describing a basic error diffusion technique. The continuous-tone input level for one pixel of the input image is shown as X. For purposes of illustration it will be assumed that the continuous-tone input level spans the range from 0 to 255. The input level X is added to a weighted error signal Xfe using a summer 10 to produce a modified input level Xin. (The generation of the weighted error signal Xfe will be discussed shortly.) The modified input level Xin for the current pixel is quantized using a quantizer 20 to form the quantized output level Xo. For a binary error diffusion algorithm, the quantizer 20 will output a 0 for any input level below a threshold (typically code value 128), and a 255 for any input level above the threshold. A summer 30 receives the modified input level Xin and the quantized output level Xo and produces an error signal Xerr representing the error introduced by the quantization process. The error signal Xerr is multiplied by a series of error feedback weights using a weighted error generator 40 to produce the weighted error signal Xfe, which is added to the continuous-tone input levels of nearby pixels which have yet to be processed using summer 10. The propagation of the errors made during the quantization process to the nearby pixels insures that the arithmetic mean of the pixel values is preserved over a local image region.\nFIG. 2 shows a typical set of error feedback weights used in the prior art. The error for the current pixel, located at (column,row)=(i,j) of the image, is weighted by the error feedback weights, and added to nearby pixels yet to be processed. In this example, there are four error feedback weights having the same value of ¼, with the sum of the error feedback weights being 1. This ensures that the arithmetic mean of the image is preserved.\nWhen using the standard error diffusion algorithm to process a color image, a technique common in the prior art is to apply the error diffusion algorithm independently to each color channel. See, for example, U.S. Pat. No. 5,757,517 to Couwenhoven, et al. This arrangement is shown for processing an image with cyan (C), magenta (M), and yellow (Y) color channels in FIG. 3. This arrangement is successful at preserving the arithmetic mean of the image data in each color channel, resulting in the correct output color, but suffers from the fact that the placement of the output dots in each color channel is decorrelated, and overprints can unnecessarily occur, resulting in an undesirable noisy, grainy appearance to the printed image. To illustrate this, consider that a 16×16 pixel image containing a 30% tint of uniform CMY code values is to be halftoned to 2 levels (on, off) and printed on an inkjet printer. As shown in FIG. 4, the CMY color channels, when processed independently with the error diffusion algorithm, produce dot patterns that typically resemble the cyan, magenta, and yellow patterns shown as patterns 50, 60, and 70, respectively. In these patterns, the black pixels indicate locations where an ink dot will be printed, and white locations indicate white paper. When these three patterns are printed on top of each other on the page, an output pattern 80 will result. (Due to the black and white reproduction of output pattern 80 in this document, only the luminance of the pattern is shown). Each pixel in the output pattern 80 can be one of eight colors: white; C; M; Y; red (C+Y); green (C+Y); blue (C+M); or black (C+M+Y). The luminance modulation between the darker pixels (R, G, B, or black) and the lighter pixels (white, C, M, or Y) produces an undesirable grainy appearance to the human eye. Interestingly, this is not necessary, since the original input image was a 30% tint of CMY, therefore it should have been possible to produce the output image using non-overlapping patterns such that only white, C, M, or Y are produced at each pixel. The luminance modulation between these colors would be much less, providing a more pleasing result.\nError diffusion methods that attempt to provide correlation between the color channels are known in the prior art. These algorithms are sometimes called “vector error diffusion” algorithms. Chapter 16 of “Digital Color Halftoning”, by H. Kang (SPIE Optical Engineering Press, 1999) describes several vector error diffusion techniques. U.S. Pat. No. 5,375,002 to Kim, et al., discloses a color error diffusion method in which an error look-up table is used to help determine the color of an output pixel. U.S. Pat. No. 5,565,994 to Esbach discloses an error diffusion method in which output signals from one color channel are used to determine threshold values for other color channels to provide some decorrelation between the color channels. U.S. Pat. No. 6,637,851 to Van de Velde, et al., discloses an error diffusion method in which a luminance signal is computed and separately error diffused as an additional channel. The luminance information is then used to guide the selection of the output pixel colors. Typically, the prior art methods are computationally complex, requiring substantial additional processing power or memory.\nIn U.S. Patent Application Publication No. 2008/0055649, Spaulding et al. have disclosed an error diffusion algorithm incorporating error signal offsets which provides significant image quality improvements with only a modest increase in the computational complexity relative to prior art algorithms. (This algorithm will be discussed in more detail later.) However, it has been found that even this algorithm can run 20-40% slower than algorithms that do not use error signal offsets.\nTherefore, there is a need for a color error diffusion method which provides for high quality printed images with reduced graininess, and can be implemented efficiently without requiring substantial additional computing power or memory."}
{"text": "Catheterization and interventional procedures, such as angioplasty or stenting, are generally performed by inserting a hollow needle through a patient's skin and tissue into the vascular system. A guide wire may be advanced through the needle and into the patient's blood vessel accessed by the needle. The needle is then removed, enabling an introducer sheath to be advanced over the guide wire into the vessel, e.g., in conjunction with or subsequent to a dilator.\nA catheter or other device may then be advanced through a lumen of the introducer sheath and over the guide wire into a position for performing a medical procedure. Thus, the introducer sheath may facilitate introducing various devices into the vessel, while minimizing trauma to the vessel wall and/or minimizing blood loss during a procedure.\nUpon completing the procedure, the devices and introducer sheath are removed, leaving a puncture site in the vessel wall. Traditionally, external pressure is applied to the puncture site until clotting and wound sealing occur. However, the patient must often remain bedridden for a substantial period of time after clotting to ensure closure of the wound. This procedure, however, may be time consuming and expensive, requiring as much as an hour of a care provider's time. It is also uncomfortable for the patient and requires that the patient remain in the operating room, catheter lab, or holding area. In addition, a risk of hematoma exists from bleeding before hemostasis occurs."}
{"text": "There are various types of electronic devices that comprise thin-film magnetic components, such as magnetic-random-access-memories and magnetic recording heads. The thin-film magnetic components in these devices are often processed by standard photolithography and etching techniques during fabrication. For electronic devices having magnetic thin-film components and stacks of other thin-film components of different natures or chemical properties, it becomes difficult to efficiently and successfully processing the thin-films using standard lithography and etching techniques. It becomes even more difficult to use standard photolithography and etching techniques to process the thin-films when the thin-films are to be defined into features with characteristic dimensions matching today's technology nodes, such as 130 nm or less.\nTaking magnetic-random-access memories (MRAM or MRAM cell) as an example, MRAMs are a new non-volatile memory technology and have been drown great attention in both scientific research laboratories and industries. Their advantageous properties over existing memory technologies for storing digital signals have proved MRAMs to become a promising mainstream memory technology in the recent future.\nA MRAM cell uses a magnetic tunnel junction (MTJ) as a storage element; and the MTJ comprises of two magnetic layers separated by a thin (such as 1 nm) insulating layer. One of the two magnetic layers, which is referred to as a reference layer, is characterized by fixed magnetization. The other magnetic layer, which is referred to as a storage layer, is characterized by variable magnetization orientation.\nThe two magnetic layers of the MTJ are often based on 3d metals (such as Fe, Co, Ni) and 3d metal alloys. The insulating layer laminated between the two magnetic layers in the MTJ often comprises of alumina (Al2O3) or magnesium oxide (MgO) although many other oxide/nitride materials could in principle be used. In one example, one of the two magnetic layers of the MTJ is made of a synthetic antiferromagnet that involves ultra thin layers of Co-alloys and ruthenium (Ru). One or both of the magnetic layers in the MTJ can be coupled with an anti-ferromagnetic layer, which can be a Mn alloy, such as FeMn, PtMn, and IrMn.\nA MRAM cell may comprise additional functional elements, such as, buffer layers to promote adhesion and texture, capping layers prevent corrosion or materials inter-diffusion, electrical contact layers, thermal barriers, and spin polarizing layers. Because of different desired functions of different elements, a MRAM cell may comprise of various materials, some of which can be uncommon materials such as NiFe, CuN, NiFeCr, Pt, GeTeSb, and BiTe, as well as usual semiconductor materials such as Ti, TiN, TiW, TiWN, W, Ta, Cu, and CoSiN.\nThe combination of many materials with very different chemical natures makes it difficult if not impossible to etch using existing semiconductor processing techniques, especially when it is to be patterned into small individual elements (“cells”) at features sizes matching today's technology nodes (130, 90, 65 nm going down to 45, 32, 22 nm). In addition, it is always desired to preserve the chemical/crystalline nature of the tunnel barrier/storage and reference layers interfaces so as to achieve desired electrical characteristics of the MRAM cell. In particular it is always desired to avoid disturbing the magnetic properties of the reference and storage layers. It is also desired to avoid variation of the critical dimension of the MRAM cell. It is also desired to avoid damages of the tunnel barrier layer by means of atomic diffusion of metallic species and/or modification of oxygen content. It is further desired to avoid electrical shorting of the tunnel barrier layer by metallic sidewalls re-depositions during fabrications.\nAmongst the existing etch techniques that are commonly used in the semiconductor and thin film industry, wet etch is unsuitable for processing MRAM features, especially those features with critical dimensions. Ion beam etching (IBE) is unsuitable either for processing MRAM features due to the following reasons. An IBE etch is often driven by high energy ions that sputter off the target material. Heavy sidewall re-deposition occurs as the sputter species being non-volatile by nature. Although etching at a grazing incidence may reduce sidewall re-deposition, such grazing incidence is primarily practical for isolated devices (such as recording heads) but not for dense MRAM cells, such as an array of MRAM cells, which is necessary for practical memory applications. Moreover, IBE etch may result in etched sidewalls being slopped due to non-isotropic etch, which in turn, causes severe critical dimension gain at the tunnel junctions of MRAM cells.\nReactive ion etching, which is capable of achieving features of critical dimensions in MRAM cells and clean vertical sidewalls, is however difficult to implement due to the multiple and often chemically incompatible elements in MTJs of MRAM cells. This arises from the fact that some of MTJs are comprised of highly non volatile elements (such as Pt and Co). Some of MTJs are comprised of highly volatile elements (such as Ge and Te); while some of MTJs are comprised of elements that are highly sensitive to corrosion (such as elements Fe and Ni). Some MTJs are comprised of elements prone to solid state diffusion (such as elements Mn, Cu, and Sb). It is therefore very difficult to find an appropriate combination of chemistry and etching parameters (such as temperature and power) to achieve a proper etch.\nMTJs of MRAM cells patterning is currently performed either by means of IBE in the data storage market, where structures are isolated by nature and grazing incidences can be used, or by means of RIE in MRAM applications, wherein large densities are required. It is believed however that as the feature size is decreased and the complexity of MTJ stacks are increased, existing etching techniques will become more and more difficult to implement in processing MRAM cells.\nIn view of the foregoing, it is desired for a method of processing thin-film components in electronic devices."}
{"text": "Heat shock protein 90 (Hsp90), one of the most abundant proteins expressed in cells, regulates cellular homeostasis by chaperoning protein folding and trafficking. Hsp90 is also highly upregulated in response to stress. The N-terminal domain of Hsp90 includes an ATP binding site, and ATPase activity is necessary for all of its cellular functions. To date, over 200 Hsp90 “client” proteins have been identified and many of these are involved in signal transduction.\nHsp90 has been implicated in diseases such as cancer, and its expression is up-regulated during oncogenesis. High expression is associated with poor prognosis for cancers, such as breast cancer, prostate cancer, non-small cell lung cancer, colorectal cancer and melanoma.\nProteins such as glucose-regulated protein 94 (GRP94) and tumor necrosis factor receptor-associated protein (TRAP1) share homology with Hsp90 and both proteins also possess ATPase activity. While inhibitors of Hsp90 have been shown to have antiproliferative and antitumor activities, current Hsp90 inhibitors may bind nonspecifically to GRP94 and TRAP1. There is a need for compounds that selectively bind to Hsp90."}
{"text": "Generally, this kind of speed-changing device is so constructed that the control lever is operated to move forward to thereby pull the wire so that the four member linkage mechanism can be transformed against the return spring provided at the four member linkage mechanism, whereby the chain guide can move forward. The control lever is also operable to return so as to allow the return spring to return so that the chain guide can also return.\nHowever, in this construction, when the control lever is operated while applying thereto a rotational resistance which overcomes the return spring, it is necessary to keep the control lever in a predetermined functional position. Hence, especially when the control lever is operating against the return spring, it is necessary to apply to the lever a force which overcomes the force of the return spring and the rotational resistance, so that a significant amount of force is required to operate the control lever and a light touch speed-change operation cannot be performed.\nWith respect to the above problem, it has been proposed that one of the four members which constitute the linkage mechanism at the derailleur be provided with a positioning plate which has a number of recesses therein corresponding to the number of speed-changing stages. Another member moving relative to the member which is provided with the positioning plate is provided with a ball engageable with the recesses. When the above two members move relative to each other, the ball travels with respect to the recesses to thereby engage with a selected one of the recesses of a selected speed-changing position, whereby the speed-changing positions of the chain guide are determined and these positions can be maintained.\nHowever, in the above construction, because the ball is resiliently engaged with the recesses for maintaining the positions of the chain guide the return spring can not be used to cause the chain guide to return through the returning function of the return spring. As a result, it is necessary to utilize two wires of a general pull-type wire or to utilize a push-pull type wire for forcibly reciprocating the chain guide.\nAccording to the above construction, since the return spring is not used, operation of the lever can be performed only by a force with which the ball can disengage from the recesses, whereby operation of the lever can be performed lightly to that extent. However, this construction cannot utilize a pull-type wire as the control wire thereby causing a problem that the speed-change device is expensive to manufacture."}
{"text": "1. Field of the Invention\nThe invention relates to an organic electroluminescence element (hereinafter, referred to as an organic EL element), a display device including the organic EL element, and a method of manufacturing the organic EL element.\n2. Background Art\nAn organic EL element has a configuration in which an organic luminescent layer showing at least an electroluminescence phenomenon is interposed between an electrode as an anode and an electrode as a cathode. The organic EL element is a spontaneous-luminescent element in which the organic luminescent layer emits light when a voltage is applied between the electrodes, holes and electrons are injected into the organic luminescent layer, and the holes and the electrons are recombined.\nTo improve the luminescent efficiency, a hole-injection layer, a hole-carrying layer, or an electron-block layer is suitably selected and formed between the anode and the organic luminescent layer, and a hole-block layer, an electron-carrying layer, an electron-injection layer, or the like is suitably selected and formed between the organic luminescent layer and the cathode. Further, the organic luminescent layer is combined with the hole-injection layer, the hole-carrying layer, the electron-block layer, the hole-block layer, the electron-carrying layer, the electron-injection layer, and the like, and the combined layers are referred to as a luminescent medium layer.\nEach layer of the luminescent medium layers is formed of an organic material or an inorganic material. The organic material is classified into a low-molecular material and a high-molecular material.\nExamples of the low-molecular material include copper phthalocyanine (CuPc) for the hole-injection layer, N,N′-dephenyl-N,N′-bis(3-methyl phenyl)-1,1′-biphenyl-4,4′diamine (TPD) for the hole-carrying layer, tris(8-quinolinol)aluminum (Alq3) for the organic luminescent layer, 2-(4-biphenylyl)-5-(4-tert-butyl-phenyl)-1,3,4,-oxazole (PBD) for the electron-carrying layer, and LiF for the electron-injection layer.\nEach layer of the luminescent medium layer formed of the low-molecular material is formed so as to have a thickness of about 0.1 nm to about 200 nm mainly by a vacuum deposition method such as a resistance heating method or a vacuum drying method (dry process) such as a sputtering method.\nThe kinds of low-molecular materials are abundant, and an improvement in the luminescent efficiency, luminescent luminance, lifetime, or the like is expected by combinations thereof.\nExamples of the high-molecular material include, for the organic luminescent layer, a material in which a low-molecular luminescent dye is solved in high molecules such as polystyrene, polymethylmethacrylate, or polyvinyl carbazole, a high-molecule fluorescent substance such as a polyphenylene vinylene derivative (hereinafter, referred to as PPV) or a polyalkyl fluorene derivative (hereinafter, referred to as PAF), and a rare-earth metal-based high-molecule phosphor.\nThe high-molecule material is generally solved or dispersed in a solvent and is film-formed with a thickness of about 1 nm to about 100 nm by a wet method (wet process) such as a coating or printing method.\nWhen the wet method is used, it is possible to obtain the advantages that a film can be formed in the atmosphere, equipment is cheap, an increase in size is easy, and a film can be efficiently formed in a short time, compared to a case where a vacuum drying method such as a vacuum deposition method is used.\nAn organic thin film formed of the high-molecule material has the advantages that crystallization or agglomeration rarely occurs, defects such as a short-circuit or a dark spot can be prevented from occurring due to coating of a pin hole or a foreign substance of another layer.\nOn the other hand, examples of the inorganic material include, for the carrier-carrying layer, alkali metal elements such as Li, Na, K, Rb, Ce, and Fr, alkali earth metal elements such as Mg, Ca, Sr, and Ba, lanthanoids such as La, Ce, Pr, Nd, Sm, Eu, Gd, Db, Dy, Ho, Er, Tm, Yb, and Lu, actinoids such as Th, metal elements such as Sc, Ti, V, Cr, Fe, Co, Ni, Cu, Zn, Y, Ar, Nb, Mo, Ru, Pd, Ag, Cd, Hf, Ta, W, Re, Os, Ir, Pt, Au, Al, Ga, In, Sn, Tl, Pb, and Bi, and metalloid elements such as B, Si, Ge, As, Sb, and Te and include inorganic compounds such as alloys, oxides, carbides, nitrides, borides, sulfides, and halides thereof.\nSince many inorganic materials have higher adhesion or thermal stability than organic materials, suppression of an erroneous luminescent phenomenon caused by a leakage current, reduction in the generation of non-luminescent regions called dark spots, and an improvement in luminescent characteristics or lifetime are expected for the inorganic materials. Further, since the inorganic materials can be applied to large-sized displays and mass-produced products at relatively lower costs compared to the organic materials, the inorganic materials are more advantageous in terms of reduction in costs.\nThere are known configurations (see JP-A-11-307259, JP-A-5-41285, JP-A-2000-68065, JP-A-2000-215985, JP-A-2006-114521, and JP-A-2006-155978) in which an inorganic hole-injection layer formed of an inorganic material is formed between an organic luminescent layer and an anode serving as a hole-injection electrode.\nFurther, there are known configurations (see JP-A-2002-367784, JP-A-5-41285, JP-A-2000-68065, and JP-A-2006-155978) in which an inorganic electron-injection layer formed of an inorganic material is formed between the organic luminescent layer and a cathode serving as an electron-injection electrode.\nWhen a process subsequent to a process of forming the inorganic material layer includes a manufacturing process of forming a film in the atmosphere by a wet method, a problem may arise in that the light luminescent efficiency, the luminescent luminance, the lifetime, or the like may deteriorate since an undesirable phenomenon such as a change in the physical properties of the inorganic material occurs due to a liquid used in the process or the oxygen or moisture of the atmosphere.\nFor example, molybdenum oxide is known as a useful material which is an example of a hole-carrying material or an electron-injection material, since film formation is easy, a hole-injection function from a hole-injection electrode is high, a function of reliably carrying holes is excellent, and stability is high.\nThe molybdenum oxide is classified broadly into molybdenum trioxide (MoO3) and molybdenum dioxide (MoO2). Since the transmittance of molybdenum trioxide is high and molybdenum dioxide is low at the time of forming a film, molybdenum trioxide is generally used.\nHowever, molybdenum trioxide is slightly soluble in water. Therefore, since molybdenum trioxide is formed and then reacts to moisture, the physical properties of molybdenum trioxide may easily vary. On the contrary, since molybdenum dioxide or most of the other inorganic compounds is insoluble in water, the physical properties of molybdenum dioxide or most of the other inorganic compounds scarcely vary.\nIn particular, when a process of forming a luminescent medium layer adjacent to molybdenum oxide includes a manufacturing process, such as transportation or film formation, performed in the atmosphere, the film may deteriorate due to a deterioration causing factor such as moisture. For this reason, a problem may arise in that display characteristics such as luminescent efficiency, luminescent luminance, or lifetime may deteriorate in some cases.\nThat is, when the luminescent medium layer is laminated only by a vacuum drying method, the amount and influence of deterioration causing factor adsorbed on the surface of molybdenum trioxide is small. However, when the manufacturing process of forming the film in the atmosphere is included, the display characteristics may considerably deteriorate.\nFurther, when a water-based solvent, an alcohol-based solvent, a ketone-based solvent, a carboxylic acid-based solvent, a nitryl-based solvent, or an ester-based solvent is used in the manufacturing process of forming the film in the atmosphere as in a wet method, a problem may arise in that particularly the luminescent efficiency and the luminescent luminance may deteriorate since molybdenum trioxide is solved and the physical properties and the film thickness are changed.\nAs a method of preventing the influence of the deterioration causing factor, there is a known method of preventing the deterioration causing factor by containing another metal compound having a strong environmental tolerance or a strong solvent tolerance in the metal compound used to form one layer of the luminescent medium layer and using the metal compound as a mixture film.    Patent document 1: JP-A-H11-307259    Patent document 2: JP-A-2002-367784    Patent document 3: JP-A-H05-41285    Patent document 4: JP-A-2000-68065    Patent document 5: JP-A-2000-215985    Patent document 6: JP-A-2006-114521    Patent document 7: JP-A-2006-155978"}
{"text": "This invention relates to tetracyclic antidepressants, dibenzocycloheptenes with a nuclearly fused cyclopropyl moiety.\nDostert and Kyburz (U.S. Pat. No. 3,803,234) disclose N-(aminoalkyl or heterocyclic aminoalkyl) dibenzocycloheptene imine derivatives.\nJudd (U.S. Pat. No. 3,349,128) discloses dibenzocycloheptenes having an aminoalkylene ether attached to the imino nitrogen but without the nuclear fused cyclopropyl moiety.\nThe following patents relate to compounds with the same basic carbon skeleton but different substituents in the 6-position: Coyne (U.S. Pat. No. 3,574,199), aminoalkyl or aminoalkylidene; Coyne (U.S. Pat. No. 3,547,933 and U.S. Pat. No. 3,658,908), hydroxy or lower alkanoyloxy and various aminoalkyl groups; Remy (U.S. Pat. No. 3,475,438), N-alkyl piperidine; Uyeda (Belgian Pat. No. 805,433 and French Pat. No. 2,201,089), alkylated imine.\nMental illness encompasses both psychoses and neuroses. Symptoms requiring treatment include depression, anxiety, agitation, and hallucinations. Among the drugs used particularly for treatment of both reactive and endogenous depressions are monoamine oxidase (MAO) inhibitors, such as iproniazide, tranylcypromine, nialamide, phenelzine, and pargyline, and the non-MAO-inhibiting tricyclic aromatic dibenzazepines, such as imipramine, and dibenzocycloheptenes such as amitriptyline.\nAll of these drugs have adverse side effects that limit their usefulness. MAO inhibitors may benefit milder forms of depression, but the risk of serious toxic effects is a strong argument against their use. They may cause liver damage and acute hypertension, especially if given in conjunction with cheese, bananas, or other amine-containing foods. The MAO inhibitors may also cause tremors, insomnia, hyperhydrosis, agitation, hypermanic behavior, confusion, hallucinations, convulsions and orthostatic hypotension. They frequently cause dizziness, vertigo, headache, inhibition of ejaculation, difficulty in urination, weakness, fatigue, dry mouth, constipation and blurred vision.\nImipramine may cause blurred vision, dryness of mouth, constipation, urinary retention, orthostatic hypotension, respiration depression, myocardial infarction, and congestive heart failure. Similar difficulties are experienced with amitriptyline.\nThere is a continuing need for psychotherapeutic agents that have fewer side effects than the drugs in use today; also for psychotherapeutic agents that have different modes of action than presently used agents, since none of these is completely effective.\nThe present invention results from efforts to develop new, safe, and effective psychotherapeutic compounds with minimal side effects."}
{"text": "1. Technical Field of the Invention\nThis invention relates generally to wireless communication and more particularly to transceivers used to support wireless communications in unlicensed spectra.\n2. Description of Related Art\nCommunication systems are known to support wireless and wireline communications between wireless and/or wireline communication devices. Such communication systems range from national and/or international cellular telephone systems to the Internet to point-to-point in-home wireless networks to radio frequency identification (RFID) systems. Each type of communication system is constructed, and hence operates, in accordance with one or more communication standards. For instance, wireless communication systems may operate in accordance with one or more standards including, but not limited to, RFID, IEEE 802.11, Bluetooth, advanced mobile phone services (AMPS), digital AMPS, global system for mobile communications (GSM), code division multiple access (CDMA), local multi-point distribution systems (LMDS), multi-channel-multi-point distribution systems (MMDS), and/or variations thereof.\nDepending on the type of wireless communication system, a wireless communication device, such as a cellular telephone, two-way radio, personal digital assistant (PDA), personal computer (PC), laptop computer, home entertainment equipment, RFID reader, RFID tag, et cetera communicates directly or indirectly with other wireless communication devices. For direct communications (also known as point-to-point communications), the participating wireless communication devices tune their receivers and transmitters to the same channel or channels (e.g., one of the plurality of radio frequency (RF) carriers of the wireless communication system) and communicate over that channel(s). For indirect wireless communications, each wireless communication device communicates directly with an associated base station (e.g., for cellular services) and/or an associated access point (e.g., for an in-home or in-building wireless network) via an assigned channel. To complete a communication connection between the wireless communication devices, the associated base stations and/or associated access points communicate with each other directly, via a system controller, via the public switch telephone network, via the Internet, and/or via some other wide area network.\nFor each wireless communication device to participate in wireless communications, it includes a built-in radio transceiver (i.e., receiver and transmitter) or is coupled to an associated radio transceiver (e.g., a station for in-home and/or in-building wireless communication networks, RF modem, etc.). As is known, the receiver is coupled to the antenna and includes a low noise amplifier, one or more intermediate frequency stages, a filtering stage, and a data recovery stage. The low noise amplifier receives inbound RF signals via the antenna and amplifies then. The one or more intermediate frequency stages mix the amplified RF signals with one or more local oscillations to convert the amplified RF signal into baseband signals or intermediate frequency (IF) signals. The filtering stage filters the baseband signals or the IF signals to attenuate unwanted out of band signals to produce filtered signals. The data recovery stage recovers raw data from the filtered signals in accordance with the particular wireless communication standard.\nAs is also known, the transmitter includes a data modulation stage, one or more intermediate frequency stages, and a power amplifier. The data modulation stage converts raw data into baseband signals in accordance with a particular wireless communication standard. The one or more intermediate frequency stages mix the baseband signals with one or more local oscillations to produce RF signals. The power amplifier amplifies the RF signals prior to transmission via an antenna.\nCurrently, wireless communications occur within licensed or unlicensed frequency spectrums. For example, wireless local area network (WLAN) communications occur within the unlicensed Industrial, Scientific, and Medical (ISM) frequency spectrum of 900 MHz, 2.4 GHz, and 5 GHz. While the ISM frequency spectrum is unlicensed there are restrictions on power, modulation techniques, and antenna gain. Another unlicensed frequency spectrum is the V-band of 55-64 GHz.\nDifferent radio networks sometimes share the same spectrum. For example, Bluetooth transceivers and 802.11g transceivers may both be present in a single area using the 2.4 GHz band. In the V-band, devices using Wireless HD (WiHD) and devices using the Next Generation Microwave System (NGMS) may be present in a single area. Transmissions by one device can cause interference with other devices that use the same frequency band with the same area.\nOther disadvantages of conventional approaches will be evident to one skilled in the art when presented the disclosure that follows."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for characterizing a specimen of semiconductor material and, more particularly, to a method and apparatus determining parameters such as the doping concentration profile of a specimen of semiconductor material.\n2. Description of Prior Art\nVarious methods for measurements of semiconductor surface and semiconductor/insulator interface parameters have been described. They can be classified into destructive and non-destructive. Examples of destructive methods are the secondary ion mass spectrometry and the four-point sensor method.\nNon-destructive methods can be further classified into those requiring post-processing and those that don't (post-processing is defined as additional processing steps beyond those that define properties of semiconductor surface or semiconductor/insulator interface that are the goal of characterization). An example of a method that requires post-processing are traditional Metal-Insulator-Semiconductor Volt-Farad (MIS CV) measurements. They require presence of an insulator over semiconductor surface and deposition of metal layer over the insulator. Additional steps are costly, time-consuming and they can change the properties that were the goal of measurements.\nMethods that do not require post-processing can be subdivided into contact and non-contact methods. Contact methods are defined as those that involve bringing the specimen into physical contact with materials or chemicals other than the ambient environment or materials and equipment that are used to handle specimens (semiconductor wafers) during normal processing (cleaning solutions, wafer chucks, etc.). An example of such methods in the area of semiconductor surface/interface characterization is a Mercury Sensor. Since such methods require direct contact of a foreign material with the front side of the wafer there is a risk of contaminating the wafer, therefore they are not used with production wafers.\nFinally, several non-contact methods for semiconductor surface/interface characterization are known in the field.\nThe apparatus described in U.S. Pat. No. 5,233,291 uses traditional CV methodology for measurement of the properties of semiconductor surface and semiconductor/insulator interface. The difference between this apparatus and MIS CV measurement systems is in the use of precisely controlled air gap as the insulating layer between the metal electrode and the semiconductor. Wafer proximity is determined by detecting the energy losses of the laser beam that undergoes total internal reflection on the surface of the transparent conductive electrode caused by interaction of the sample with the evanescent beam. The electrode is suspended at submicron distance from the sample, this distance as well as tip and tilt being monitored by three capacitive sensors and controlled by three piezo actuators in real time during the measurements.\nIn order to calculate semiconductor parameters using CV methods, it is necessary to know the width of the depletion layer in the semiconductor W.sub.d. In this method W.sub.d is calculated from the capacitance of the depletion layer, which in turn is calculated from the total capacitance between the metal electrode and the semiconductor substrate. EQU C.sub.tot.sup.-1 =C.sub.air.sup.-1 +C.sub.SiO2.sup.-1 +C.sub.si.sup.-1,\nthus it is important to know C.sub.air at all times during the measurement, which means either measuring the air gap or holding it fixed with substantial precision during the measurements of electrical parameters of the sample, both tasks being technically complicated, considering the submicron width of the air gap.\nFurthermore, this method requires knowledge of variations of C.sub.si which is obtained from measurements of C.sub.tot. The presence of two additive components in the equation above, C.sub.air.sup.-1 and C.sub.SiO2.sup.-1 means that the sensitivity of the capacitance meter has to be lower than desirable for measurements of variations in C.sub.si.sup.-1.\nAdditionally, the method described in U.S. Pat. No. 5,233,291 is not well suited for measurements of samples with poorly passivated semiconductor/insulator interface, e.g. non-passivated wafers of epitaxially grown films, because of the presence of slow surface states which may recharge during the measurements and thus affect both the measurements of C.sub.si.sup.-1 and the readings of the three capacitive gap sensors.\nThe surface photovoltage method described in U.S. Pat. No. 4,827,212 to E. Kamienicki, which patent is incorporated herein by reference makes use of modulated light whose wavelength corresponds to an energy greater than the band gap of the semiconductor and whose modulation frequency is greater than the reciprocal lifetime of minority carriers and whose intensity is low enough so that the modulation of the semiconductor surface potential is small compared to the surface potential. The photovoltage generated by the light is proportional to the width of the depletion region near the semiconductor surface. The photovoltage can be measured using capacitive coupling of a conductive electrode to the surface of the sample. The advantage of this method over regular CV methods is in the fact that it is not required to measure separate contributions of the insulating gap capacitance or the oxide capacitance which makes the method less complicated and less sensitive to errors in sensor positioning. However, this method is not a true non-contact method because it contacts the specimen with an insulating film.\nIn U.S. Pat. No. 5,453,703 to W. C. Goldfarb, which patent is incorporated herein by reference, a method and apparatus are disclosed for determining the minority carrier surface recombination lifetime constant (t.sub.s) of a specimen of semiconductor material. The specimen is positioned between a pair of electrodes, the specimen being disposed on one of the electrodes and being spaced from the other electrode. A signal is provided corresponding to the capacitance between the specimen and electrode spaced from the specimen. A region of the surface of the specimen is illuminated with a beam of light of predetermined wavelengths and which is intensity modulated at a predetermined frequency and varying in intensity over a predetermined range. A fixed bias voltage V.sub.g applied between the pair of electrodes, the fixed bias voltage being of a value such that the semiconductor surface is in a state of depletion or inversion. A signal is provided representing the ac photocurrent induced at the region of the specimen illuminated by the light beam. The intensity of the light beam and frequency of modulation of the light beam are selected such that the ac photocurrent is nearly proportional to the intensity of the light beam and reciprocally proportional to the frequency of modulation of the light beam. A signal is provided corresponding to the illumination intensity of the beam of light. The surface minority carrier recombination time constant (t.sub.s) is then determined using the ac photocurrent capacitance and illumination intensity information.\nIt is an object of this invention to provide a new and improved method and apparatus for characterizing a semiconductor.\nIt is another object of this invention to provide a new and improved method and apparatus for determining the doping concentration profile and average doping concentration of a specimen of semiconductor material.\nIt is another object of this invention to provide a non-contact method and apparatus for determining the doping concentration profile and average doping concentration of a specimen of semiconductor material.\nIt is still another object of this invention to provide a method and apparatus as described above in which photovoltage (or photocurrent) is measured using capacitive coupling of a sensor electrode to the surface of the specimen and which does not require precise control or direct monitoring of the gap between the sensor and the specimen.\nIt is a further object of this invention to provide a method and apparatus as described above which is particularly suitable for characterization of specimens with high density of slow surface states, e.g. non-passivated wafers.\nIt is still a further object of this invention to provide a new and improved sensor assembly for use in non-contact measuring of semiconductor surface photovoltage or photocurrent.\nIt is another object of this invention to provide a new and improved capacitive pickup type sensor assembly.\nIt is still another object of this invention to provide a non-contact method of measuring semiconductor surface properties."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of fabricating a semiconductor acceleration sensor consisting of a supporting part, a beam part and a weight part made of semiconductors used for sensing an acceleration by detecting the deflection of the beam part accompanying the displacement of the weight part.\n2. The Description of the Prior Art\nA sensor for detecting an acceleration by finding the amount of deflection, due to the acceleration, of a beam part having a weight, from the change in the resistance value or the like has already been known widely, and miniaturization of the acceleration sensor through formation of the sensor using silicon as the semiconductor is being put into practice. As an example of such a semiconductor acceleration sensor, one may mention the device disclosed in a paper entitled \"A Batch-Fabricated Silicon Accelerometer\", IEEE Electron Devices, Vol. ED-26, No.12, p.1911, December 1979. According to this paper, an acceleration sensor consisting of a supporting part 31, a beam part 32 and a weight part 33 as shown in FIG. 1 can be formed by subjecting a single sheet of silicon substrate to anisotropic etching.\nIn such a sensor, the shape of the weight part and the beam part directly affects the sensitivity of the sensor, and in particular a high dimensional accuracy is demanded for the thickness of the beam part.\nThe fabrication process of a semiconductor acceleration sensor according to anisotropic etching is shown in FIGS. 2A to 2D as a first example of the prior art. In this method, first, using thermal oxide films 42 and 43 on a silicon substrate 41 as the etching masks (FIGS. 2A and 2B), the region between a supporting part 44 and a weight part 45 is etched from the rear surface by anisotropic etching (FIG. 2C), then a gap part 46 surrounding the periphery of the weight part is formed by giving the anisotropic etching also from the front surface (FIG. 2D).\nAs techniques of accomplishing a high precision working of the beam thickness, methods of using a silicon substrate consisting of two crystal layers with different conductivity types, and removing the crystal layer of one conductivity type and using the crystal layer of the other conductivity type as the beam part will be described as a second and third examples of the prior art in the following.\nFIGS. 3A to 3F show the fabrication process of a second example of the prior art disclosed in Japanese Laid-Open Patent Publication No. 63-292071. According to this method, first, an n-type silicon substrate 52 is formed on a p-type silicon substrate 51, and a thermal oxide film 53 is formed on both surface of the laminated substrate (FIG. 3A). The thickness of the n-type silicon substrate 52 is the same thickness of a desired beam thickness using the thermal oxide film 53 as the mask, the n-type silicon substrate 52 and a part of the p-type silicon substrate 51 are removed in a V-shaped groove 54 to form a gap part (FIG. 3B). Next, a metallic film 55 for establishing electrical conduction with the n-type silicon substrate 52 is formed on the n-type silicon substrate (FIG. 3C). Next, the p-type silicon part between the supporting part 56 and the weight part 58 is removed by using an electrolytic etching, or electrochemical etching from the p-type silicon substrate 51 side (FIG. 3E). When the etching is reached to the metallic film 55, electrochemical etching is stopped automatically. Thus the thickness of the n-type silicon substrate 52 becomes the beam thickness.\nIn FIGS. 4A to 4G, the fabrication process of a third example of the prior art is shown. First, an n-type silicon substrate 62 is formed on a p-type silicon substrate 61, and a thermal oxide film 63 is formed on both sides of the laminated substrate (FIG. 4A). After removing the p-type silicon part in the region between the supporting part and the weight part by applying an electrochemical etching from the rear surface side using the thermal oxide film as the etching mask (FIG. 4C), the entire rear surface is covered with a protective film 64 of gold or the like (FIG. 4D), and a gap part is formed by carrying out an anisotropic etching from the front surface side (FIG. 4F). Finally, the sensor is completed by removing the gold protective film 64 (FIG. 4G).\nAlthough several fabrication methods of the prior art semiconductor acceleration sensor have been proposed, these prior art fabrication methods have the following drawbacks. To begin with, in the first example of the prior art fabrication method, the thickness of the beam part varies with the temperature of etchant at the time of etching, the etching time and the stirring condition of the etchant. Furthermore, the etching takes a long time because of the large amount to be etched as much as 200 to 400 .mu.m, so that high accuracy control of the beam thickness is not easy to accomplish.\nIn the second example of the prior art, as shown in FIG. 3D, since the bottom portion of the V-shaped groove is dug into the p-type silicon substrate, the electrochemical etching in the step shown in FIG. 3E does not proceed beyond the bottom portion of the V-shaped groove. Thus a desired beam thickness cannot be achieved.\nFurthermore, in the third example of the prior art, the covering property of the rear surface protective film becomes a matter of concern, where pin holes are generated also in the beam part due to the pin holes in the protective film, giving rise to a drawback that about 40% of units produced are defective."}
{"text": "Field of the Invention\nOne or more embodiments of the invention are related to the field of containers. More particularly, but not by way of limitation, one or more embodiments of the invention relate to a cup lid with integrated container that enables simultaneous or intermittent access of the contents of the container and attached cup without disengagement of the cup lid from the cup. Additionally, an independent drop-in container may reside within the cup lid cavity such that after partially consuming the contents of the independent drop-in container it may be resealed with an additional lid and removed from the lid cavity and transported to another location such as a car or home.\nDescription of the Related Art\nStandard cup lids are simple covers that do not include an integrated container. Rather, known lids cover the contents of a cup which forms a closed container in combination with the cup itself. Known containers that couple with cups include food containers that fit onto the top of yogurt cups for example. Known containers have to be removed from the yogurt cup and then flipped over and opened before the contents of the container and cup may be accessed. It is generally not possible to access the contents of the cup while also accessing the contents of the container without first disengaging the container from the cup. Additionally, food containers that attach to yogurt cups in an upside-down position have a limited food-volume capacity. In such cases, as the yogurt example shows, the food-container walls narrow as they proceed upward toward the bottom of the upside down container. Other known devices having a container or shelf combined with a lid have limitations which makes these devices impractical to use. One category of devices includes a container combined with a cup, but utilizes a hole in the middle of the lid. This makes it impossible to store relatively circular items, i.e., non-ring or non-annular items having no central hole, in the container, such as hamburgers, cookies or muffins for example. Another category of device includes a container combined with a lid, but does not allow for simultaneous access of the contents of the cup and the container at the same time, and does not allow for the container to be resealed or a drop-in container to be removed from the container. Other devices that include drop-in functionality require removal of the container before accessing the contents of the cup. Yet these devices do not contemplate a drop-in container that is configured to fit into the arm rest of a movie theater seat. Other devices have relatively small peel containers for pills such as mints and are not suitable for larger food items. Another category of devices utilizes dividers in the cup with access on each side of the cup. None of the known devices enable a container to be disengaged from the lid of the cup while retaining the lower lid on the cup. No known devices have a non-permanent or male/female bottom oriented coupling system for coupling a container with the lid. Furthermore, there are no known rotational covers that enable or disable access to the liquid and/or solid in the cup as desired by the user.\nKnown containers that couple with bottles include gift containers that fit onto the top of bottles for example. It is generally not possible to access the contents of the bottles while also accessing the contents of the gift containers without disengaging the gift container from the bottle and then disengaging the lid of the bottle.\nThus simultaneous or intermittent access of the contents of known cups or bottles and of the contents of an attached container is not possible. This makes for difficult drinking/eating of coffee, soda, snacks, popcorn, etc., in malls, fast food restaurants, theaters, amusement parks, sports stadiums or in any other venue. For example, this makes it difficult to eat and drink food in a theater or stadium with one cup-holder per seat.\nFor at least the limitations described above there is a need for a cup lid with integrated container."}
{"text": "Several 2,4-diaminopteridine derivatives being substituted in the 6-position and/or the 7-position of the pteridine ring (according to standard atom numbering for said ring) are known in the art, e.g. from various sources of literature including Swiss Patent No. 231,852; British Patent No. 763,044; U.S. Pat. No. 2,512,572; U.S. Pat. No. 2,581,889; U.S. Pat. No. 2,665,275; U.S. Pat. No. 2,667,486; U.S. Pat. No. 2,940,972; U.S. Pat. No. 3,081,230 and U.S. Pat. No. 5,047,405. Some of these substituted 2,4-diaminopteridine derivatives were disclosed in relationship with various medical uses, such as bacterial growth inhibitors, antineoplastic agents, anti-schistosomiasis activity, coronary dilating activity, diuretic and hypotensive activity, and anti-amnesic activity. In particular, U.S. Pat. No. 2,940,972 and EP-A-362,645 disclose very specific 2,4-diaminopteridine derivatives being substituted by piperidinyl, morpholinyl or pyrrolidinyl in the 7-position of the pteridine ring.\nSpecific 2-aminopteridine derivatives wherein the 4-position of the pteridine ring is substituted with an alkoxy group and the 6-position is also substituted are also known in the art, although without any medical utility. For instance, U.S. Pat. No. 2,740,784 discloses such derivatives wherein the 6-substituent is an acetal; J. Chem. Soc. (1957) 2146–2158 discloses 2-dimethylamino-4-ethoxy-6-phenylpteridine as a compound with a melting point of 200° C.; Arm. Khim. J. discloses 2-amino-4-ethoxy-6,7-diphenylpteridine; Helv. Chem. Acta (1992) 75:2317–2326 discloses 2-amino-4-pentoxy-6-methylthiopteridine. Furthermore, International Patent Application published as WO 01/19825 discloses 2-phenylamino and 2-phenylsulfoxide pteridines wherein the 7-position is substituted with an amide as being useful inhibitors of cell cycle regulatory kinases.\nNevertheless, there still is a need in the art for specific and highly therapeutically active compounds, such as, but not limited to, drugs for treating immune and autoimmune disorders, organ and cells transplant rejections, cell proliferative disorders, cardiovascular disorders, disorders of the central nervous system and viral diseases. In particular, there is a need in the art to provide immunosuppressive compounds or antineoplastic drugs or anti-viral drugs which are active in a minor dose in order to replace existing drugs having significant side effects and to decrease treatment costs.\nCurrently used immunosuppressive drugs include antiproliferative agents, such as methotrexate (a 2,4-diaminopteridine derivative disclosed by U.S. Pat. No. 2,512,572), azathioprine, and cyclophosphamide. Since these drugs affect mitosis and cell division, they have severe toxic effects on normal cells with high turn-over rate such as bone marrow cells and the gastrointestinal tract lining. Accordingly, marrow depression and liver damage are common side effects of these antiproliferative drugs.\nAnti-inflammatory compounds used to induce immunosuppression include adrenocortical steroids such as dexamethasone and prednisolone. The common side effects observed with the use of these compounds are frequent infections, abnormal metabolism, hypertension, and diabetes.\nOther immunosuppressive compounds currently used to inhibit lymphocyte activation and subsequent proliferation include cyclosporine, tacrolimus and rapamycin. Cyclosporine and its relatives are among the most commonly used immunosuppressant drugs. Cyclosporine is typically used for preventing or treating organ rejection in kidney, liver, heart, pancreas, bone marrow, and heart-lung transplants, as well as for the treatment of autoimmune and inflammatory diseases such as Crohn's disease, aplastic anemia, multiple-sclerosis, myasthenia gravis, uveitis, biliary cirrhosis, etc. However, cyclosporines suffer from a small therapeutic dose window and severe toxic effects including nephrotoxicity, hepatotoxicity, hypertension, hirsutism, cancer, and neurotoxicity.\nAdditionally, monoclonal antibodies with immunosuppressant properties, such as OKT3, have been used to prevent and/or treat graft rejection. Introduction of such monoclonal antibodies into a patient, as with many biological materials, induces several side-effects, such as dyspnea. Within the context of many life-threatening diseases, organ transplantation is considered a standard treatment and, in many cases, the only alternative to death. The immune response to foreign cell surface antigens on the graft, encoded by the major histo-compatibility complex (hereinafter referred as MHC) and present on all cells, generally precludes successful transplantation of tissues and organs unless the transplant tissues come from a compatible donor and the normal immune response is suppressed. Other than identical twins, the best compatibility and thus, long term rates of engraftment, are achieved using MHC identical sibling donors or MHC identical unrelated cadaver donors. However, such ideal matches are difficult to achieve. Further, with the increasing need of donor organs an increasing shortage of transplanted organs currently exists. Accordingly, xenotransplantation has emerged as an area of intensive study, but faces many hurdles with regard to rejection within the recipient organism.\nThe host response to an organ allograft involves a complex series of cellular interactions among T and B lymphocytes as well as macrophages or dendritic cells that recognize and are activated by foreign antigen. Co-stimulatory factors, primarily cytokines, and specific cell-cell interactions, provided by activated accessory cells such as macrophages or dendritic cells are essential for T-cell proliferation. These macrophages and dendritic cells either directly adhere to T-cells through specific adhesion proteins or secrete cytokines that stimulate T-cells, such as IL-12 and IL-15. Accessory cell-derived co-stimulatory signals stimulate activation of interleukin-2 (IL-2) gene transcription and expression of high affinity IL-2 receptors in T-cells. IL-2 is secreted by T lymphocytes upon antigen stimulation and is required for normal immune responsiveness. IL-2 stimulates lymphoid cells to proliferate and differentiate by binding to IL-2 specific cell surface receptors (IL-2R). IL-2 also initiates helper T-cell activation of cytotoxic T-cells and stimulates secretion of interferon-γ which in turn activates cytodestructive properties of macrophages. Furthermore, IFN-γ and IL-4 are also important activators of MHC class II expression in the transplanted organ, thereby further expanding the rejection cascade by enhancing the immunogenicity of the grafted organ The current model of a T-cell mediated response suggests that T-cells are primed in the T-cell zone of secondary lymphoid organs, primarily by dendritic cells. The initial interaction requires cell to cell contact between antigen-loaded MHC molecules on antigen-presenting cells (hereinafter referred as APC) and the T-cell receptor/CD3 complex on T-cells. Engagement of the TCR/CD3 complex induces CD154 expression predominantly on CD4 T-cells that in turn activate the APC through CD40 engagement, leading to improved antigen presentation. This is caused partly by upregulation of CD80 and CD86 expression on the APC, both of which are ligands for the important CD28 co-stimulatory molecule on T-cells. However, engagement of CD40 also leads to prolonged surface expression of MHC-antigen complexes, expression of ligands for 4-1 BB and OX-40 (potent co-stimulatory molecules expressed on activated T-cells). Furthermore, CD40 engagement leads to secretion of various cytokines (e.g., IL-12, IL-15, TNF-α, IL-1, IL-6, and IL-8) and chemokines, all of which have important effects on both APC and T-cell activation and maturation. Similar mechanisms are involved in the development of auto-immune disease, such as type I diabetes. In humans and non-obese diabetic mice, insulin-dependent diabetes mellitus results from a spontaneous T-cell dependent auto-immune destruction of insulin-producing pancreatic .beta. cells that intensifies with age. The process is preceded by infiltration of the islets with mononuclear cells (insulitis), primarily composed of T lymphocytes. A delicate balance between auto-aggressive T-cells and suppressor-type immune phenomena determines whether expression of auto-immunity is limited to insulitis or not. Therapeutic strategies that target T-cells have been successful in preventing further progress of the auto-immune disease. These include neonatal thymectomy, administration of cyclosporine, and infusion of anti-pan T-cell, anti-CD4, or anti-CD25 (IL-2R) monoclonal antibodies. The aim of all rejection prevention and auto-immunity reversal strategies is to suppress the patient's immune reactivity to the antigenic tissue or agent, with a minimum of morbidity and mortality. Accordingly, a number of drugs are currently being used or investigated for their immunosuppressive properties. As discussed above, the most commonly used immunosuppressant is cyclosporine, which however has numerous side effects. Accordingly, in view of the relatively few choices for agents effective at immunosuppression with low toxicity profiles and manageable side effects, there exists a need in the art for identification of alternative immunosuppressive agents and for agents acting as complement to calcineurin inhibition.\nThe metastasis of cancer cells represents the primary source of clinical morbidity and mortality in the large majority of solid tumors. Metastasis of cancer cells may result from the entry of tumor cells into either lymphatic or blood vessels. Invasion of lymphatic vessels results in metastasis to regional draining lymph nodes. From the lymph nodes, melanoma cells for example tend to metastasize to the lung, liver, and brain. For several solid tumors, including melanoma, the absence or the presence of lymph nodes metastasis is the best predictor of patient survival. Presently, to our knowledge, no treatment is capable of preventing or significantly reducing metastasis. Hence, there is a need in the art for compounds having such anti-metastasis effect for a suitable treatment of cancer patients.\nIn the field of allergy, IgE is well known for inducing allergy mainly by stimulating mast cells to release histamine. Also, asthma, being characterized by inflammation of airway and bronchospasm, is mainly induced by Th2 cytokines such as IL-5, IL-10 or IL-13. Therefore there is a need in the art for compounds that efficiently inhibit the release of these Th2 cytokines.\nThere is also a need in the art to improve therapeutic efficiency by providing pharmaceutical compositions or combined preparations exhibiting a synergistic effect as a result of combining two or more immunosuppressant drugs, or antineoplastic drugs or anti-viral drugs or anti-histamine drugs.\nMeeting these various needs in the art constitutes the main goal of the present invention."}
{"text": "Tinnitus is the perception of sound in the absence of corresponding external sounds. Tinnitus may be caused by injury, infection or the repeated bombast of loud sound, and can appear in one or both ears. Although known for its high-pitched ringing, tinnitus is an internal noise that varies in its pitch and frequency. The sound perceived may range from a quiet background noise to a signal loud enough to drown out all outside sounds.\nTinnitus affects approximately 50 million Americans or about 15% of the general population. There is no effective treatment to ameliorate tinnitus. The prior art has included devices and methods that attempt to mask tinnitus.\nGeneral Hearing and several others manufacture an earpiece that can generate sound. The masking methods work well for people who suffer mild forms of tinnitus but don't work for people who suffer from loud tinnitus because to mask the tinnitus, the external sound has to be louder than the tinnitus.\nSome treatments, such as drugs, surgery, psychotherapy and masking, exist, but none are consistently effective and may have significant side effects. Here we propose to use sound to suppress tinnitus. Different from previous masking methods, which use an external sound that is typically louder than the tinnitus and has similar pitch quality to the tinnitus, we propose to use optimized sounds that may be softer and have different pitch qualities than the tinnitus. Another significant difference between previous masking devices and the present suppressing method is that the masking effect is instantaneous, while the present suppressing method usually takes time to develop and decay. The sound can be delivered either acoustically via an audio device (e.g., MP3 players, iPod® devices, or other ear devices) or electrically via a temporary or permanent implant.\nIn particular, various medications have been used to treat tinnitus, which are aimed at the disease or its co-morbidities. Tinnitus is often associated with anxiety and depression. Whether the condition causes psychological distress or whether tinnitus is more prevalent in people with psychological disorders is under debate. It has been suggested that tricyclic antidepressants treat the underlying psychological problems rather than the tinnitus directly. Some studies have demonstrated significant decrease in tinnitus intensity with the administration of intravenous (IV) lidocaine but have not reproduced the effect with oral analogues such as tocainide. IV administration of lidocaine is not a practical method due to its short half-life and serious side effects. Other agents, such as carbamazepine, alprazolam, baclofen, betahistine, and cinnarizine have also been studied as possible tinnitus treatments, with inconsistent results. Selective serotonin reuptake inhibitors have shown no benefit or a slight benefit in placebo-controlled studies.\nTinnitus masking devices of the prior art deliver a continuous low-level noise or tone to the ear that results in camouflaging of the tinnitus. Use of the device has been shown to reduce tinnitus severity. The Neuromonics tinnitus treatment is a new method of acoustic stimulation in combination with a structured program of counselling. The device output is a broadband stimulus that is constructed for each individual's hearing profile. Three clinical trials reported general improvement in quality of life of patients using this treatment. Hearing aids also grant partial or total relief from tinnitus by amplifying external sound stimuli thereby decreasing the awareness of the presence of tinnitus. Hearing aids are usually tried as first line treatment in patients with tinnitus in the presence of hearing loss. Hearing aids can also be used in combination with sound generators that increase the rate of success.\nIn the end of 1980s the neurophysiological model of tinnitus was proposed which then elicited the development of a new therapeutic approach called tinnitus retraining therapy (TRT). TRT consists of counselling along with delivery of a low-level, constant white noise to the affected ear(s) of a patient. The sound usually results in habituation of the patient's auditory system to the tinnitus thereby decreasing the patient's awareness of their condition. The rate of success varies between different patients. TRT requires one-to-two years of therapy for effective management of tinnitus.\nCochlear implantation has been found to be effective in reducing tinnitus in 28% to 79% of patients who have associated bilateral severe sensorineural hearing loss in previous studies. Exacerbation of tinnitus is reported in 9% of patients after cochlear implantation. Success with cochlear implants in decreasing tinnitus has led to other forms of electric stimulation, including transtympanic stimulation of the cochlea, vestibulocochlear nerve stimulation, and transcranial magnetic stimulation of the auditory cortex. Recent studies have shown that 1 out of 3 cochlear implant users and 5 out of 11 (45%) transtympanic subjects showed substantial or total tinnitus relief while stimulated by high rate pulse trains. More than that, the external perception was either presented subthreshold or diminished to below threshold within a short period of time. This is the ideal situation for tinnitus suffers since both the tinnitus and external stimulus are not perceived as a result. In a review of the literature, unilateral cochlear implant use was suggested to be associated with decline of contralateral tinnitus in 67% of patients. The effect of a cochlear implant is demonstrated to be due to masking of tinnitus by surrounding sounds that are newly perceived or from electrical stimulation of the auditory nerve. The mechanisms underlying the complexity of tinnitus are unclear, and a cure for this condition remains elusive. The prevailing theories generally point to either a central or peripheral source of subjective tinnitus, defined as sound perceivable only to the patient without a clear organic cause. The reality most likely involves a combination of both, with peripheral injuries and pathology causing functional and architectural changes in the central nervous system.\nUnited States Patent Application Publication 2007/0203536 (Hochmair et al.) describes cochlear implant devices as well as other implantable devices usable to deliver tinnitus treatment. The entire disclosure of United States Patent Application Publication 2007/0203536 is expressly incorporated herein by reference.\nU.S. Pat. No. 7,347,827 (Choy) describes monofrequency tinnitus patient treatment apparatus and process wherein phase cancellation effects are purported to be achieved by utilizing an externally generated sound which is subjectively selected by the monofrequency tinnitus patient to match in both tone and loudness his or her tinnitus sound. This subjectively selected externally generated sound wave which matches in tone and loudness the patient's tinnitus sound, is either (i) sequentially phase shifted through a plurality of phase shift sequence steps totaling at least 180 degrees or (ii) alternatively is directly phase shifted in essentially a single step motion into a 180 degree, out-of-phase reciprocal, canceling relationship with the patient determined tinnitus tone. The sequential steps of the phase shifted tone or the directly phase shifted tone are applied to the tinnitus patient to effect cancellation or diminishment of the patient's tinnitus. The entire disclosure of U.S. Pat. No. 7,347,827 (Choy) is expressly incorporated herein by reference.\nPhase cancellation techniques, such as that described in U.S. Pat. No. 7,347,827 (Choy) may not be optimally effective in treating tinnitus. Most tinnitus is of high frequency (>1000 Hz) although some low-frequency tinnitus cases do exist. However, even if a person's tinnitus is low frequency, the phase canceling techniques may be ineffective or less than optimal because the perceived tinnitus sound is a neural impulse not an acoustic sound and cannot be physically or externally cancelled by phase inversion. Physiologically all sounds will be half-wave rectified in the cochlea and become digital-like pulses (called action potentials) in the nervous system. An out-of-phase external sound will add, rather than cancel, the internally-generated tinnitus.\nTinnitus occurs in the setting of sensorineural hearing loss in the majority of patients, thus the postulate that tinnitus originates peripherally in the cochlea. Noise damage causes molecular changes to structural proteins in stereocilia and the cuticular plate. Cytoplasmic calcium levels increase dramatically in response to sound, potentially disrupting normal hair cell function. Progressive insult results in complete destruction of hair cells in certain regions of the basilar membrane. Aberrant auditory signals occur around areas of hair cell loss, an edge effect that results in the perception of sound. Furthermore, auditory nerve fibers are spontaneously active during quiet, resulting in neurotransmitter release. Loss of the spontaneous activity can lead to abnormal central auditory activity, which could be perceived as sound. Lack of sound input and edge effects as a cause of tinnitus could explain the reduction of tinnitus commonly seen following cochlear implantation. Tinnitus suppression by electrical stimulation in hearing and deaf patients has been shown in a relatively recent study.\nOn the other hand, not all patients with hearing loss experience tinnitus, and not all tinnitus patients have hearing loss. Sectioning of the auditory nerve can result in tinnitus, suggesting a central etiology. Positron emission tomography (PET) imaging of patients with tinnitus shows unilateral neural activity in Brodmann areas 21,22, and possibly 42. These effects were decreased with lidocaine administration, a known suppressor of tinnitus. Sounds caused bilateral activity regardless of which ear was stimulated. Patients who experienced tinnitus in response to facial movement showed unilateral changes in cerebral blood flow. Another study using single photon emission computed tomography (SPECT) imaging during tinnitus showed a 16% increase in signal in the right auditory cortex and a 5% increase in the left, whereas noise caused bilateral change. The effect was eliminated by lidocaine administration. This difference in the activation pattern supports a theory of central generation of tinnitus.\nThis interpretation is consistent with animal models of tinnitus showing changes in spontaneous activity of the dorsal cochlear nucleus (DCN) following noise or cisplatin exposure. This hyperactivity persists following ablation of the damaged cochlea, suggesting that while a peripheral injury may stimulate a change in the DCN, the central hyperactivity is independent of peripheral input.\nTinnitus has been separated in three categories depending on severity and if hearing loss is present. Category 0 is characterized by the tinnitus having a low impact on the person's life. Categories 1 and 2 are used to describe tinnitus with a high impact on life with Category 2 indicating the presence of subjective hearing loss that accompanies the tinnitus. For our studies this convention was followed but a second parameter based on loudness was defined. The loudness parameter was determined by each subject as Low, Moderate or High. On a 10-point scale with 1 being the lowest sound (threshold) and 10 being the upper limit of loudness, Low is 0 to 3, Moderate is 3 to 6, and High is 6 to 10. Careful measures were taken to ensure each subject understands they are ranking loudness and not annoyance.\nFIG. 1 of this patent application is a graphic representation of tinnitus severity classification and the typical tinnitus patient population. Patients within category 0 are least likely to seek tinnitus treatment. Category 1 and 2 patients with low levels of loudness are likely to be helped by TRT. It has been shown that TRT helps in relaxation for a majority of patients (72.5%) but benefits a much smaller percentage with respect to their ability to work (25.5%) and sleep (47%). At present, category 1 and 2 patients with high levels of loudness are often left without effective treatment.\nThere remains a need in the art for the development of new devices and methods for the treatment of tinnitus."}
{"text": "The new Sutera cultivar is a product of a planned breeding program conducted by the inventor, Gabriel Danziger, in Moshav Mishmar Hashiva, Israel. The objective of the breeding program was to produce plants of Sutera cordata with compact habits and large pink flowers. The seed parent is the unpatented, proprietary seedling variety referred to as Sutera cordata cv.5-259. The pollen parent is unknown as ‘DANCOP27’ is a result of a plannned open pollination in which the pollen parent is impossible to distinguish.\nThe new variety was discovered in May 2006 by the inventor, in a group of seedlings resulting from that crossing, in Moshav Mishmar Hashiva, Israel.\nAsexual reproduction of the new cultivar ‘DANCOP27’ by vegetative cuttings was first performed at a commercial greenhouse in Moshav Mishmar Hashiva, Israel in June 2006 and has shown that the unique features of this cultivar are stable and reproduced true to type in successive generations."}
{"text": "Drugs are commonly administered in solid form through pills or capsules that can be orally taken. However, many biological drugs cannot be administered this way because of degradation in the gastrointestinal tract and quick elimination by the liver. Another common technique for administration of drugs in liquid form is through injection using a metal hypodermic needle that can cause significant pain and discomfort to patients. For example, U.S. Pat. No. 5,279,585 discloses an injection device is provided for injecting fluids such as insulin within body tissue. The device includes a housing, a piston rod movably mounted within the housing, and a dose setting mechanism for controlling the movement of the piston rod with respect to a fluid-containing cartridge removably positioned within the housing.\nSimilarly, a number of physical and chemical techniques including electroporation, laser ablation, ultrasound, thermal, iontophoresis, and chemical enhancers have been explored to develop painless transdermal drug delivery techniques. It was found that its very difficult for the molecules with a molecular weight higher than 500 or diameter larger than 1 nm to penetrate normal human skin. Further studies showed that the key barrier for transdermal delivery of substances is the stratum corneum layer, the outer layer of skin, that is about 4-30 micron thick. Invasive methods to overcome this skin barrier have been used in practice, such as intradermal (ID), intramuscular (IM), or subcutaneous (SC) injection using standard hypodermic needles and syringes. These methods cause pain and require a skilled professional. In addition, they may cause needle injuries. Similarly, current methods of extracting biologic fluids such as blood from patients suffer from the same disadvantages.\nIn order to improve the skin permeability of the therapeutic agents and other active ingredients, microneedles have been recently developed to disrupt the stratum corneum and facilitate the delivery of the active agents and ingredients to the epidermis. These active substances can then diffuse through the rest of the epidermis to the dermis and be absorbed by blood vessels and lymphatics there. The substance absorbed can then get into the circulation system. Thus, both topical and systemic-level delivery of drugs is possible. Since there are no nerves and blood vessels in the stratum corneum and epidermis, this is a minimally invasive, painless and blood-free method of drug delivery. An additional advantage of this method, when engineered for topical delivery of vaccines, can lead to an enhanced inoculation effect because the epidermis is rich in antigen presenting cells and is a desired target for vaccine delivery."}
{"text": "The present invention relates to an electrode for a living body, which is suitably used as the electrode of a defibrillator, an electrocardiograph and the like.\nHitherto, in the case of carrying out the catheter-inspection of the heart, a catheter is inserted into a blood vessel of the patient and after introducing a contrast media into the heart via the blood vessel, X-ray photography is carried out.\nHereupon, there are cases where the patient is exited and causes cardiac convulsion for longer than 15 to 30 sec during the inspection. In the case where the cardiac convulsion continues for a long time, the stagnation of the blood occurs, thereby causing the damage of the cerebral cells, etc. Accordingly, a high voltage shock of 3000 to 5000 V is given to the heart of the patient while using a defibrillator ordinarily.\nThe conventional defibrillator has a shape of a box provided with a handle on the upper part thereof, and the electric shock is given by pressing the electrode part which is disposed on the lower surface of the box on the region of the heart of the patient.\nHowever, in the case of using the conventional defibrillator, it is necessary to use the defibrillator after taking out it on each time when the cardiac convulsion is caused on the patient.\nNamely, it is very much troublesome to use the conventional defibrillator and particularly, there is a fear of too late to meet an emergent requirement.\nMoreover, there are cases where an electrode has been preliminarily adhered onto the surface of the body of the patient due to the necessity of taking electrocardiogram during X-ray photographing. However, since the conventional electrode is made of a metal and does not transmit X-ray, such an electrode has been an obstacle to the X-ray photography.\nAs a result of the present inventors' studies for solving the above-mentioned problems, it has been found that the whole electrode can be made to be transparent to X-ray by making the electrode substrate of the porous material mainly consisting of granular or fibrous carbon or of the flexible and porous material mainly consisting of fibrous carbon, and on the basis of the above-mentioned finding, the present inventors have attained the present invention.\nNamely, the object of the present invention is to provide an electrode for a living body which has been constituted so that the whole electrode is transparent to X-ray and does not hinder the X-ray inspection by making the electrode substrate of the porous material composed of granular or fibrous carbon, or of the flexible and porous material composed of fibrous carbon."}
{"text": "A company route is defined by a list of records that include airports, procedures, airways, and/or waypoints satisfying regulatory authorities regarding the flight of an aircraft. Essentially, a company route is a description of the flight path to be followed by an aircraft as it flies between origin and destination airports. This flight path is designed specifically to fit the requirements of airline operators who operate between origin and destination points.\nA company route is typically designed on the ground by navigation personnel, who receive the requirements from the airline operators. The navigation personnel use a list of waypoints between origin and destination airports along the flight path as inputs so as to produce the company route. The design of a company route is a mostly manual process and it is certainly time-consuming. As a result, the process is error prone. Also, the design process is textually oriented, which makes it difficult for the navigation personnel to visualize the actual company route, which is made even more difficult because, during this design process, the designer has no idea about the terrain and obstacles which may present along the designed company route.\nAs can be seen, the step by step process that is necessary to create a company route is complex. The source airport, the departure, runway transitions, common transitions, en-route transitions, waypoint fixes, airways, arrival transitions, approach transitions, and the actual approach all/few need to be selected and specified based on defined criteria and stringed together in a way that forms the record of the company route. At each selection point, the navigation personnel use multiple sources (charts, Aeronautical Information Packages (AIPs), etc.) to select the desired points of interest. Because so many sources are needed to provide the information required for designing a company route, there is a distinct possibility of making mistakes, and the designer will not have any leverage to view the sources together. Also, during this whole process, the designer has no idea about the terrain and obstacles which may be present along the designed route. Moreover, text based route design provides little means to select an optimal path.\nTherefore, there is a need for an interactive and graphical navigation system that interactively provides navigation personnel with a graphical view of a company route as it is being designed, along with additional information such as information about terrain, obstacles, airports, airways, etc. This system can be used by navigation personnel on the ground and can be arranged to empower them to be more aware of the choices of the route at a particular fix (a real time view of the route on a world map backdrop if a choice is selected), to compare different choices so that they can make better decisions quicker, to be aware of terrain and obstacle information, and to verify the company route by viewing a virtual fly through of the route.\nThe relevant choices (e.g., fixes, airways, procedures, etc.) relative to a particular fix can be shown graphically, and the user can make selections based on the graphics so as to construct the company route. This company route can then be verified by virtually flying through the coded route with underlaid terrain and obstacles."}
{"text": "Medical syringes are used to inject a fluid or medication inside a body. A plunger is pushed along the inside of the syringe barrel, which then moves the fluid through a needle inside the body. FIG. 1 shows a typical medical syringe barrel 100 with plunger 110. Plunger 110 is typically a rubber tip with for example three ribs.\nLubrication of the inside surface 120 of the syringe barrel 100 is crucial to reduce the friction between the plunger 110 and the inner surface of the syringe barrel and ensure effective injection of fluid 130.\nSyringe barrels are normally lubricated by spraying the inside of the syringe barrel using a spraying nozzle, which is fixed in space in the middle of the cylindrical barrel, and which in turn is moved up and down past the sprayer while it is active.\nThis method of lubrication generally works well, but is subject to irregularities, which can arise when the sprayer does not spray uniformly in all directions or may not be centered in the middle of the syringe barrel. The present invention advances the technology of syringe barrels by addressing at least some of the current problems with lubrication of these syringe barrels."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for controlling equipment in a semiconductor manufacturing system, and more particularly, to a method for automatically controlling process conditions that reflect actual operational capabilities of units within the equipment, through an unit monitoring module, to maximize the productivity of the system.\n2. Description of the Related Art\nGenerally, the fabrication of semiconductor devices involves highly precise processes that require finely tuned precision equipment. Several pieces of precision equipment are typically employed in sequence and arranged on a semiconductor processing line. The operation of each piece of precision equipment on the line is closely monitored by operators to maintain and enhance the efficiency of the processing line.\nAs shown in FIG. 1, conventional fabrication equipment 3 are disposed on a conventional processing line. When a lot (not shown) of workpieces, such as wafers, are introduced into the equipment 3, the equipment 3 performs a fabrication process on the lot. The equipment 3 is connected on-line to a host computer 1 through an equipment server 4. An operator interface (O/I) 2, for example an operator interface personal computer (O/I PC), is also connected on-line to the host computer 1. Through the O/I 2, an operator informs the host computer 1 that a process using the equipment 3 is about to commence. The operator inputs basic manufacture data into the host computer 1, e.g., the identification number (ID) of the lot to be processed in the equipment 3 and the ID of the equipment 3 for performing the process with the lot.\nThen, based on the input basic manufacture data, the host computer 1 searches its data base for the process condition data to be applied to the process on the lot. The host computer 1 immediately downloads these process condition data, including preset process settings, to the equipment. Process settings may include, for example, a desired process time duration or a desired process temperature.\nThereafter, the operator checks the process condition data and inputs a process commencing command or a process terminating command. The lot is then automatically and rapidly routed into and out of the equipment 3. In this manner, the process equipment 3 performs on the workpieces of each lot based on the process settings received.\nThe equipment 3 performs the processes on the lots by operating units or components of the equipment 3, e.g., loading and unloading ports (not shown), and chambers 5. While the equipment 3 performs the process on the lots, the operator closely monitors the operations of the units of the equipment 3 continuously. If it is recognized that a particular unit, e.g., a particular chamber 5, of the equipment 3 is in an non-operational state, because, for example, the particular unit is undergoing preventative maintenance or it is out of order, then the whole equipment is determined to be in the non-operational state. Consequently, the equipment is turned off to prevent production problems associated with introducing lots into non-operational chambers.\nIn the processing line, a transferring apparatus, e.g., an auto guided vehicle (AGV) 6, is included to automatically transfer lots among the equipment 3 and between bays of the fabrication facility. The AGV 6 is controlled by commands communicated from the host computer 1 to an on-line transfer apparatus server 7 and then transmitted by the transfer apparatus server to the AGV 6 via a communications link. The host computer 1 continuously communicates operational commands to the AGV 6 so that the AGV 6 can rapidly respond to perform the lot transferring function.\nHowever, such a conventional controlling system suffers from several problems. First, even when only some of the units of the equipment 3 are non-operational and other units are operational, the operator considers the whole equipment 3 to be in the non-operational state. As a result, the equipment 3 is not fully utilized and the productivity of the processing line is not efficiently controlled.\nSecondly, since it is very difficult for the operator to recognize the actual operating states of units in real time, equipment 3 with non-operational units may be considered to be in the operational state for some period of time. During this time delay, the operator may introduce lots into the equipment with process condition data designed for normal operations. In this case, an excessive number of lots may be introduced into the equipment compared to the number the equipment can actually handle. This over-introduction of lots may result in the accumulation of lots in the equipment. When lots accumulate in the equipment, the efficiency of the entire processing line is reduced.\nThirdly, even when the operator recognizes that some units are non-operational, the AGV is operated without any intervention of the operator, so that the over-introduction and accumulation of lots in the equipment may still occur, thus increasing the chances of over-introduction."}
{"text": "The present invention relates to a non-impact printing device for recording graphic symbols on ordinary paper by means of selective emission of ink particles by an ejector.\nDevices of the aforesaid type are known wherein drops of liquid ink are emitted selectively through a nozzle by an ejector device actuated through the medium of a piezoelectric crystal and wherein to every electrical stress applied to the crystal there correspond a compression of the volume of the chamber containing the ink and a corresponding emission of drops through the nozzle.\nDevices are moreover known wherein the emission of drops is caused by applying a potential difference between the conductive liquid ink and the mouth of the nozzle, as in the U.S. Pat. No. 1,958,406.\nOne of the problems which arise with the aforesaid devices known in the art is that the liquid ink encrusts the nozzle and ends by blocking it in the course of time. The choice of special water-based inks does not permit the problem of encrustation to be solved completely."}
{"text": "Unless otherwise indicated, the foregoing is not admitted to be prior art to the claims recited herein and should not be construed as such.\nSwitching amplifiers, such as Class-D audio amplifiers, offer high efficiency operation due to the use of switching output stages. Accordingly class-D amplifiers are commonly used for audio presentation in portable device applications such as MP3 players, smartphones, computer tablets, and other battery operated devices, where power efficiency is an important consideration. DC offsets inherent in the amplifier circuits comprising class-D amplifier designs, however, impact their effectiveness in portable devices. DC offsets manifest themselves audibly as clicks and pops in the audio output, which can negatively impact the user experience; and DC offsets tend to increase quiescent power consumption, which can negatively impact power performance.\nTypical class-D amplifier designs use a 3-state pulse width modulation (PWM) modulator. The 3-state PWM modulator allows for filter-less designs, thus avoiding the cost of reconstruction filters in terms of silicon area on the chip. 3-state PWM modulators, however, can suffer from cross-talk between the plus and minus PWM channels. The cross-talk can arise from coupling with comparators in the PWM modulator design or between the comparators and the output stages of the PWM modulator. Cross-talk can degrade the total harmonic distortion (THD) characteristics of the class-D amplifier caused by PWM misfiring due to the cross-talk."}
{"text": "(1) Technical Field\nThe present invention relates to coin operated telephones and more particularly to coin operated relay device circuitry as used in prepay type telephone paystations.\n(2) Background Art\nQuite a number of different devices have been used to collect or refund coins deposited in telephone paystations. For example, in the current type 102B telephone paystation as manufactured by GTE Communication Systems, a single coil solenoid is used with a magnetic actuator. Depending on the direction of current supply to the solenoid coil, the magnet is either attracted or repelled from one magnetic pole to the solenoid, thereby causing either the collect or refund operation mode of the device.\nAnother type of coin relay used on the type 62 and type 82 prepay paystations manufactured by Automatic Electric Company, consisted of two coils, a permanent magnet between the coils and a centrally located armature mounted above the coils and magnet. The armature would then pivot to either side as determined by the polarity of the direct current applied to the relay coils. The position of the armature would then determine the coin collect mode or refund mode.\nBoth the single-coil solenoid device and the double-coil device require high voltages and currents to operate under control of an operator located at the telephone central office. A minimum of 48 volts and 44 milliamps of current are required to operate the device included in the current type 120B coin telephone referred to above. Another coin relay made by Northern Electric Company type P10E786 is similar in operation and construction to that utilized in the type 120B telephone noted above.\nThe automatic telephone exchange associated with coin telephones of previous design in prepay operation must of necessity be equipped with coin control repeaters, sources of positive and negative 110 volts DC for use as collect or refund battery and an interrupter that produces intermittent coin-control current for application to the telephone line. Upon lifting the handset in such an arrangement the calling party is connected to the central office and dial tone is usually heard. At this point the calling party now deposits the initial rate deposit and may dial the call in the usual manner. Upon termination of a call, the central office then sends either collect or refund battery to the telephone, depending upon whether the call was completed or was incomplete. If the call was completed, operation of the coin relay directs the coins into the cash vault included in the telephone paystation. If the call, however, was incomplete the coins are then directed into the coin return receptacle for return to the user. After the coins have been collected or refunded, the telephone is returned to its normal condition, that is to say it is ready for another call.\nOn an operator assisted calls, the initial rate amount is refunded upon connection to the operator. The toll operators then supervise collection of coins by audible signals generated by an audio oscillator located within the telephone. The operator controls the application of the coin collect or refund signals from the central office on all toll calls.\nWhen the calling party deposits coins, the coins are directed into a coin accepter-rejector mechanism. After the coins clear that mechanism, they pass through a coin chute, that contains in modern design units, three trigger switches, one for each of the three denominations of coins utilized (5 cent, 10 cent, and 25 cent). When the coins actuate the trigger switches in the coin chute, the trigger switches close a path to a totalizer. The action of the totalizer is not pertinent to the present invention. However it is noted that the first coin in its path to the coin chute cause the triggering of the coin relay and operates a set of contacts that prepare a ground path to the transmission network. When the initial rate has been deposited, the totalizer causes the two relays to operate thereby enabling the rotary dialer touch calling unit (TCU) and extending coin relay ground to the transmission network. At this time, the central office is prepared to receive address signalling from the dial or touch calling unit.\nIn existing typical paystations such as the type 120B manufactured by GTE Communication Systems, passage of a coin into the hopper forces the coin trigger downward in the hopper slot, tripping it to permit the coin operated contacts at the rear of the spring pileup to assume their rest position. The first operating set of contacts remove a short circuit from the dial pulse springs or a low resistance shunt from the oscillator of the touch calling unit connected at two fixed terminals. The other contacts connect ground to the line instead of to the relay coil. After the coin is past, the trigger is prevented from restoring fully by the tilt of the lever spring of the make contact set. If a -110 volts DC refund potential is applied to the line, the polarity of the flux generated by the core of the relay coil is such as to cause the left end of a small bar magnet to be attracted to the left extension of the core or pole piece, this tilts a nylon selector card slightly so that as it is moved down by a rocker arm, it moves a cam and therefore a nylon hopper vein to the right. Meanwhile, the downward motion of the card is imparted to a trap lever, which pivots on its pin and moves its upper edge outward. This permits the nylon coin trap pivoted above the lever, to swing down inside the hopper, so the coin deposit falls against the left side of the vein and passes into the refund chute. If a +110 volts DC collect potential is applied to the line, the flux generated at the pole piece extensions are such as to attract the right end of the selector card magnet, against the right side of the vein into the coin receptacle.\nAs the relay armature operates the selector cards of the rocker arm, it also moves the spring lift to place 1,000 Ohm resistance in series with the coil and then short-circuits the coil. The resistance thus takes the place of the coil in the path of the coin control current, stabilizing the latter to prevent damage to the ground contacts when they open and to resistance lamps in the supply circuit at the central office. At the end of its stroke, the spring lift provides substantial pull to the coin operated ground contacts, so that the tilt of the lever spring no longer interferes with the trigger, which restores completely to normal under the force of a coil spring in torsion. The dial or touch calling unit shunting springs are forced close. With the relay shorted the current induced by the collapsing field tends to retard the decay of the flux in the core and makes the armatures slow to release. As it begins to restore, the spring lift permits the coin operated springs to retrieve from their over travel position. This prevents the break (either from the dial or from the touch calling unit shunt) contacts sets from reopening and allows the make (ground) contact set to open. As the ground path is opened the relay is disconnected from the coin battery circuit, thus extinguishing the coin pilot lamp if collection of refund is being controlled through a manual switchboard and preventing reoperation of the relay if the operator or control circuit continues application of coin batteries to the line. Finally, the transfer context set is restored to remove the short circuit from the relay coil and bypass the series resistor. The upward motion of the selector card during release acts to return the coin trap to horizontal and the hopper vein to its neutral vertical position in preparation for the next deposit.\nFrom the forgoing it will be obvious that existing state-of-the-art coin collect and refund devices and the circuits associated therewith require high voltages and currents to operate and control by an operator or from central office equipment. Accordingly, it is the object of the present invention to provide a new circuit to be included in the telephone instrument for stand alone operation on a two-wire business line (not requiring operator or central office control)."}
{"text": "Following the end of the Second World War, chlorofluorocarbons (organic compounds that include carbon, chlorine and fluorine, hereinafter “CFCs”, as well as those chlorofluorocarbon compounds having at least one hydrogen present, hereinafter “HCFCs”) came into popular use as refrigerators replaced the traditional ice boxes. Refrigerants are capable of removing heat by cooling a space below the ambient temperature, and are typically used in a heat cycle, wherein the refrigerant undergoes a reversible phase change from a gas to a liquid, during which time heat is removed from a space. Thus, in addition to use in refrigerators, refrigerants also find uses in residential and business air conditioners, as well as in automobile air conditioning systems. Refrigerants have other uses as well, including use as foam blowing agents and in the pharmaceutical industry for asthma inhalers\nStudies, however, have linked CFCs to atmospheric ozone depletion. As a result, the Montreal Protocol (initially entered into force in 1989, and subsequently revised multiple times) mandated that CFCs no longer be used as refrigerants and their use was phased out during the period of 1987 to 1997. The use of chlorofluorocarbons as refrigerants has gradually been replaced by the use of fluorocarbons, including perfluorocarbons (FCs), hydrofluorocarbons (HFCs), and hydrochlorofluorocarbons (HCFCs). While manufacturing of chlorofluorocarbons (CFCs) has essentially ceased in view of the Montreal Protocol, existing supplies continue to have limited commercial value, for instance as inhalers in the health field. Because CFCs, HCFCs, FCs and HFCs are often present as azeotropic mixtures, each of the components of the mixture may have value, particularly if CFCs can be effectively separated from the other components. This would allow for compounds to be recycled, or for compounds to be more effectively destroyed or converted for disposal.\nThe fluorocarbon refrigerants currently in use (i.e., HCFCs, FCs and HFCs) exhibit chemical stability, non-flammability, chemical inertness, are safe to health, and do not deplete ozone from the atmosphere. The refrigerants have physical characteristics such as low melting points and boiling points appropriate for refrigerant use, low vapor heat capacity, low viscosity, high thermal conductivity, and good oil solubility.\nChlorofluorocarbons and fluorocarbons are conventionally referred to by a numbering convention wherein the rightmost digit denotes the number of fluorine atoms; the “tens” digit denotes one more than the number of hydrogen atoms; the “hundreds” digit (if present) denotes the number of carbon atoms less one (thus there may be no digit for this when the compound is a methyl halide). Further, although most of the refrigerant compounds are alkanes, if the compound contains any links of unsaturation, the number of double bonds can be indicated in numbering conversion as the “thousands” digit. There may also be a suffix, “a”, “b”, “c”, to indicate isomers. Thus, HFC-134a has 4 fluorine atoms, 2 hydrogen atoms and 2 carbon atoms. The “a” suffix indicate that the F atoms are not equally distributed, but rather the compound is 1,1,1,2-tetrafluoroethane. In contrast, in referring to HFC-134, the fluorine atoms are equally distributed throughout the molecule and the compound is 1,1,2,2-tetrafluoroethane. The alkane based compounds are generally referred to as CFCs, HCFCs, FCs, and HFCs.\nBecause chlorofluorocarbons often form azeotropic mixtures with fluorocarbons, effective and efficient separation is frequently difficult. Therefore, in view of the current restrictions associated with chlorofluorocarbons, it has become both necessary and desirable to develop new methods to effectively and efficiently separate and purify the individual components from an azeotropic mixture. Various techniques have been developed to facilitate the separation of azeotropic mixtures, including the use of various zeolites and molecular sieves, for instance, as is described in U.S. Pat. Nos. 4,906,796; 5,087,778; 5,160,499; 5,260,496; 5,288,930; and 5,585,529. Each of the prior art techniques, however, have not been entirely satisfactory because the components of these azeotropic mixtures frequently have similar boiling points, thereby making it difficult to obtain highly pure components upon separation. In addition, the prior art zeolite/molecular sieve processes frequently require that the zeolite material or molecular sieves be changed and/or regenerated, thus requiring that the separation process be stopped for a period of time. In an exemplary prior art procedure employing zeolites and/or molecular sieves, as described in U.S. Pat. No. 5,497,627, a single stage azeotropic separation of a contaminated chlorofluorocarbon (i.e., CFC-12 contaminated with HCFC-22 and CFC-115) is processed using water as a solvent. It appears very little refrigerant is extracted and water remains with both extractants.\nThe process and apparatus described herein provide an apparatus that effectively and efficiently separates components for ideal mixtures (i.e., simple binary mixtures) and azeotropic mixtures (e.g., mixtures that include one or more CFC compound, in addition to at least one compounds from the group of HCFCs, HFCs, and FCs)."}
{"text": "Field of the Invention\nThe present invention relates to a transmitter circuit, a semiconductor apparatus, and data transmission method.\nRelated Art\nIn a case where signals are exchanged between a plurality of semiconductor chips differing from each other in power supply voltage, the semiconductor chips must be electrically insulated from each other with insulating coupling elements in exchanging signals. Known insulating coupling elements include an AC coupling element using capacitors, coils and the like, and an optical coupling element (a photocoupler). Japanese Unexamined Patent Application Publication No. 2013-229812 discloses a semiconductor apparatus which exchanges signals using coils as insulating coupling elements, that is, what is called a micro-isolator.\nIn the disclosure of Japanese Unexamined Patent Application Publication No. 2013-229812, a pulse signal triggered by the edges of a data signal is transmitted from a transmitter circuit. Here, from the transmitter circuit, a pulse signal capable of distinguishing between the rising edge and the falling edge of the data signal is transmitted. Therefore, the data signal can be reconstructed at a receiver circuit.\nMeanwhile, Japanese Unexamined Patent Application Publication Nos. 2005-045100 and 2012-253241 and Japanese Patent No. 4750746 each disclose an electrostatic discharge protection circuit provided between the power supply and the ground. The electrostatic discharge protection circuit is mounted for protecting the internal circuitry of a semiconductor apparatus from high voltage pulses generated by electrostatic discharge. The electrostatic discharge protection circuit disclosed in Japanese Unexamined Patent Application Publication Nos. 2005-045100 and 2012-253241 turns on an NMOS transistor upon sensing a sharp increase in the power supply. The electrostatic discharge protection circuit (GGNMOS: Gate Grounded NMOS) disclosed in Japanese Patent No. 4750746 turns on a parasitic bipolar of an NMOS transistor when a power supply potential of a certain level is reached. By the foregoing operations, each electrostatic discharge protection circuit operates before the power supply potential reaches the breakdown voltage of the internal circuitry. Thus, an increase in the power supply voltage is suppressed and the internal circuitry is protected."}
{"text": "1. Field of the Invention\nThis invention relates to fluid dispensing systems which automatically dispense fluid to animals only during predetermined periods of time.\n2. Description of the Prior Art\nExperiments and controlled feeding often require maintaining animals on a fluid restriction regimen. These restrictions may require limiting the amount of fluid an animal may consume, the time period during which the animal may consume the fluid, or some combination of both of these restrictions. Controlled feeding programs in which a fluid is available to an animal at the same time every day present scheduling difficulties. Typically, glass bottles containing fluid are manually placed in the animal's cage each time access to the fluid is required. This manual procedure often produces cumbersome work schedules for experimentors and laboratory personnel, particularly when the experiments last several weeks or months, including weekends and holidays.\nU.S. Pat. No. 3,294,066 describes an animal feeding system wherein a time actuated valve causes a predetermined amount of water to enter a feeding unit where it will mix with a dry food material and the weight of such mixture will downwardly displace the feeding unit so as to expose a feeding nipple portion, thus allowing the animal to feed from the unit. This system, unlike the present invention, lacks a provision for draining or retrieving any fluid remaining when an animal feeding interval is over and it requires that a disposable feeding bag be replaced after each use. Additionally, this system differs from the present invention in that it is designed to dispense a fixed amount of fluid and is not adaptable to experiments where the watering time interval, rather than the quantity of fluid to be consumed, is to be restricted.\nU.S. Pat. No. 3,720,185 discloses an automatic animal feeding system in which a level sensing device controls the amount of liquid delivered to a chamber where it is mixed with a solid food material and then dispensed at individual feeding points. This system utilizes automatic valves to dispense the food but is designed to provide a predetermined amount of food to an animal rather than to allow the animal access for a limited period of time."}
{"text": "Mainstream clinical practice in endovascular treatment of intracranial aneurysms has changed little since the 1990's when vasooclusive coil use became widespread. Certainly, improved catheters and other auxiliary devices (e.g., stents) have helped make coiling procedures safer and/or more effective. However, the art in achieving adequate and appropriate aneurysm coil packing is best accomplished by the most highly skilled physicians.\nWhere practicable, aneurysm exclusion by cover-type devices (e.g., as described in U.S. patent application Ser. No. 12/397,123 to the assignee hereof) may be preferred. Certain other groups are attempting to shift the paradigm away from intra-aneurysm coil packing to achieve embolization via deployment of an extra-aneurysm flow disruptor/diverter stent in the parent vessel. These densely braided devices and/or multiple braid devices layered upon one another are placed in the parent vessel across the neck of an aneurysm with the intent to alter hemodynamics so as to effect embolization.\nThese WALLSTENT-like devices are best suited for placement across sidewall aneurysms. Yet, terminal aneurysms (e.g., bifurcation aneurysms) are estimated by some to make-up between about 60 and 80% of all aneurysm occurrences. By most optimistic count, only about 40% of intracranial aneurysms can be treated using the referenced stent-like devices.\nNumerous other devices have been conceived in effort to address terminal aneurysms. Complicated and/or impracticable deployment is common to many. Others simply serve as adjunctive to coils or liquid embolic agents. In these latter examples, procedures may become even more complicated and require even greater physician skill than a standard coiling procedure.\nA simpler, yet promising solution is proposed in PCT/US 2007/0076232 to Dieck, et al. A braided/mesh conical member is described for diverting blood flow from the aneurysm neck. A base of the device is set inside the aneurysm while a flow diverter portion extends into the parent vessel to direct blood flow toward adjacent side branches and away from the aneurysm. The implant may be positioned within the aneurysm as a stand-alone device or be supported by a connected stent-like body.\nU.S. Pat. Nos. 6,168,622 and 6,506,204 to Mazzochi, et al. discloses another type of braided flow disruptor set at least partly within an aneurysm. A bulbous portion is adapted to fit within the aneurysm dome and is anchored on the outside by a neck-covering flap. Given the manner in which bifurcation aneurysms often incorporate branch vessel anatomy, such a patch would often interfere with and/or “flap” free raising significant issues of potentially pathological thrombus formation within the parent vessel.\nImplants of the present invention address shortcomings of each of the above-referenced devices. As such, the subject implants (as well as their associated delivery systems) offer potential to advance the state of the art in endovascular treatment of vascular malformations, including aneurysms."}
{"text": "An integrated circuit may include metal gates with widths less than 30 nanometers. The metal gates may have work function layers and barriers of high electrical resistivity materials such as refractory metals. As a result, the electrical resistance per unit length of the metal gates may be undesirably high. Debiasing along the metals gates may occur in high frequency components and along metal gates used as local interconnects. Designing the integrated circuit to avoid metal gate debiasing may be problematic."}
{"text": "Induction furnaces having a crucible for receiving metallic material to be melted, hereinafter called a \"charge\", are well known. Such a crucible is generally called a \"cold crucible\" because it is not substantially heated by the induction device, and it is permanently cooled by a cooling device. In such furnaces, it is also well known to cause levitation of the charge by using electromagnetic confinement phenomenon, resulting in a separation of the charge to be melted from the inner crucible walls.\nIt is common practice to make the crucible from a plurality of metal segments, electrically isolated from each other to reduce electromagnetic losses in the segments. Collectively, the segments form the metallic wall of the crucible.\nIt is also well known to form such a crucible in a generally cylindrical shape. The bottom part of this shape is hemispheric or conical and includes a central hole for the release of the melted charge.\nA major preoccupation in the design and operation of such furnaces is avoidance of excessive heating of the crucible wall. The division of the crucible wall into a plurality of segments, electrically insulated from each other results in a reduction of the induction effects on the wall and limits heating of the wall. However, it is not possible to completely suppress such heating. Furthermore, the charge melting in the crucible transmits a certain amount of heat to the crucible wall. Therefore, it is necessary to substantially cool the segments forming the crucible wall. Generally, this cooling is carried out by forming the segments of the crucible wall of a metal, such as copper, which conducts heat, and by providing holes in the segments. These holes accommodate pipes extending parallel to the longitudinal axis of the crucible, and containing a cooling fluid such as water.\nGenerally, the segments forming the crucible are relatively thick because they have to be large enough to contain the diameter of the cooling pipes provided inside the segments. As the segments necessarily have a thickness greater than the penetration depth of the magnetic field generated by the induction system, it has not been possible to substantially avoid magnetic induction effects in these segments. Thus, in order to limit induction effects, the lateral dimensions of the crucible, i.e., its diameter, have been very limited. The width of the segments could not be made less than a determined value because they must include at least one longitudinal hole for a cooling pipe. Thus, design flexibility has been seriously curtailed.\nFor these reasons, known furnaces must use magnetic induction means operated at intermediate frequencies, i.e. greater than about 400 Hz. Such furnaces have other drawbacks such as high manufacturing costs, which increase with the number of segments used. Moreover, input and output pipes for the refrigerant must be provided for each segment. It is also necessary to provide the fluid connection of these pipes, as well as the electrical insulation separating the segments from each other. Such a crucible is therefore difficult to produce and its cost is relatively high.\nAnother drawback is that magnetic induction coils operating at intermediate frequency allow only relatively low currents. Therefore, the voltage across these coil terminals is relatively high, greater than 80 volts, for example. It is therefore necessary to take precautions with such systems to avoid risk of electrocution and to limit electric arcing.\nYet an additional drawback is that in order to supply an intermediate frequency electric current to such a magnetic induction system, it is necessary to provide an intermediate frequency electric supply transformer. Such a supply system is relatively expensive, particularly if it delivers high levels of electric energy. It has not been possible heretofore to build an induction furnace having a cooled wall, with a capacity to contain a volume of charge greater than about 5000 cm.sup.3. Generally, the inner diameter of such a crucible is never greater than 150 mm.\nIn such prior art furnaces, a problem arises concerning the mixing of metal melted in the crucible. This mixing is often not satisfactory and this has heretofore limited furnace capacity."}
{"text": "1. Field of the Invention\nThe present invention relates to laser oscillation apparatus. More particularly, this invention relates to new laser oscillator devices adaptable for use as laser-machining discharge excitation type laser oscillators which require large-output and high-efficiency laser light.\n2. Description of the Related Art\nTypically, in a discharge laser oscillation apparatus, the discharge is performed with respect to a laser medium in order to attain laser oscillation. When this is done, the laser medium increases in temperature. The higher the medium temperature, the less the efficiency of such laser oscillation. This phenomenon will become more severe especially in pulse laser oscillation devices.\nConventionally, in order to avoid such reduction of laser oscillation efficiency due to temperature increase of the laser medium thereby obtaining efficiency-enhanced laser oscillation, one approach is to wait for the laser medium to have lowered its temperature down to an acceptable level; another approach is to replace and cool the laser medium per se.\nAdditionally, in large-output lasers such as an excimer laser or carbon dioxide gas laser or the like, a convection device such as a fan blower or the like as well as a thermal exchanging machine are provided inside of a laser tube for forcibly causing a laser medium gas to exhibit convection to thereby achieve cooling by substitution and circulation from a laser cavity or between electrodes, which in turn lets the resulting laser output power stay higher.\nHowever, in the prior art, large-output fan blower devices have been required in order to attain the circulation or percolation of the laser medium such as a gas or the like within laser devices or laser tubes, which fan blowers might be configured from large-scale complicated electrical equipment including a motor, fan, vacuum shaft coupler and others, resulting in an increase in the size and production cost of the laser devices and laser tubes concerned. Another problem faced in the prior art is that resultant power dissipation can significantly increase undesirably."}
{"text": "Embodiments relate generally to methods and apparatus to implement compression and decompression of policies including, for example, hardware decompression of policy vectors associated with policies and represented by bit vectors.\nPolicies including conditions and instructions (or actions) are used in, for example, packet classification and packet routing. Known policies are often represented as text files or described in a markup language within a text file. Representations of such policies can be stored in a memory of a packet classification appliance or device, but policies often occupy a significant portion of the memory of such a packet classification device.\nCompression schemes for reducing the amount of memory occupied by a policy are known, but suffer several disadvantages. For example, compression of known representations of policies is difficult to implement in hardware-based packet classification devices due to memory, size, and logic constraints in hardware-based packet classification devices. Because such compression is difficult to implement in hardware-based classification device, benefits of hardware-based devices have not been realized in many packet classification devices. Thus, a need exists for improved methods and apparatus for improved policy classification in packet classification devices."}
{"text": "1. Field of the Invention\nThe present invention relates generally to mineral fiber production, and in particular to apparatus for limiting the emission of airborne byproduct material which is incidental to the manufacture of fiber products.\n2. Description of the Prior Art\nIn the art of producing mineral fiber, a spinning unit is utilized for fiberizing molten slag with air or steam. The fiber is conveyed through a blow chamber or collection chamber by down draft fans to a conveyer belt where it is removed by various methods and put into a conventional bagger. The disadvantage of this arrangement is that since it is not a closed system, the lint or airborne flywool byproduct cannot be completely contained, thereby allowing a large quantity of fugitive airborne particles to be emitted into the atmosphere and contributing to pollution of the environment. Because of its damaging ecological effect, this airborne, fugitive fly byproduct must be captured or recovered before the air stream in which it is entrained is discharged into the atmosphere.\nThe uncontrolled emission of these airborne particles into the atmosphere presents formidable problems because of the increasing emphasis on environmental quality which has led to surveillance and regulation by federal and local authorities. Federal and local legislation in this area no longer permit uncontrolled emissions from mineral fiber plants. Clean air legislation regulates the acceptable amount of particulate material and the constituency of effluent gases from such operations. Failure to abide by the regulatory acts can result in penalties as well as the imposition of permanent injunctions against such operation. Further, new operators must demonstrate in pilot operations that such emissions fall below the established minimum level before a permit for continuous full scale operation will be granted.\nIn view of the damaging effect of uncontrolled emissions on the ecology, the increasingly tighter controls imposed on such operations by federal and local authorities, and in view of the increasing demand for mineral fiber products, there is an urgent need for new and improved equipment which limit such emissions to safe levels while allowing the mineral fiber process to be operated efficiently on a large scale to meet product demand."}
{"text": "This invention relates to electronic devices capable of maximizing or minimizing a value which is a function of a number of parameters, and what is generically referred to as \"cost function\". The invention is more particularly directed to a multiple parameter control circuit capable of simultaneously adjusting controlled input variables for maximizing or minimizing a particular performance criterion, for example, in an industrial process.\nIn a typical industrial system, such as a machine, power plant, or the like, the operation generates a so-called \"cost function\", for example, power, efficiency, fuel economy, safety, etc. The cost function, in turn, is a function of a number of controlled parameters, such as fuel feed rate, combustion air pressure and temperature, exhaust damper settings, etc. Each of these parameters may be adjusted independently of the others, but the optimum value of the cost function will require optimal setting of all of the various parameters.\nCost functions, which can be optimized in terms of the setting of a multiplicity of input parameters, can also come about in computational problem solving, such as curve fitting, maximum likelihood calculations, polynomial expansions, and the like. Cost functions are also used in computer simulations, such as the determination of optimal designs in engineering, scheduling problems, many-body particle interactions, crystal growth problems, and other problems of physics and the sciences.\nAn orthogonal modulation technique for optimizing control systems has been described in U.S. Pat. No. 3,617,717. In that technique, a number of control signals result from multiplying a performance indicating signal by respective modulation signals, and the result is integrated and sampled. However, this technique is rather complex and so requires complex circuitry. In particular, a digital computer is required to carry out a large number of iterative calculations. The procedure is rather time consuming, which means that the cost function will not always converge quickly to an optimal value. That system can control only a small number of parameters. The correlation between feedback and cost function is weak, and convergence is slow."}
{"text": "1. Field of the Invention\nThis invention addresses a multi-wavelength fiber optic transmitter related to directly modulated laser sources and multifunctional filters.\n2. General Background and State of the Art\nFiber optic communication systems use a variety of transmitters to convert electrical digital bits of information into optical signals that are carried by an optical fiber to a receiver. In a directly modulated transmitter the output intensity of a laser is modulated by directly changing the injection current driving the laser. In an externally modulated transmitter, the intensity of a continous wave laser is modulated via the use of a modulator, which changes the intensity of the laser light. Directly modulated semiconductor lasers are typically compact, integrable, and have large responses to modulation. They are comparatively inexpensive than externally modulated transmitters, which require an intensity modulator following the laser source. However, directly modulated lasers may suffer from a drawback; namely, their outputs may be highly chirped. As a result, directly modulated lasers are normally used for short reach applications because the inherent chirp of the laser causes the transmitted pulses to be distorted after propagation in dispersive fiber. For longer reach applications, external modulation is used. However, external modulation requires a costly modulator that consumes power, introduces loss, and takes up board space."}
{"text": "The printed circuit board (PCB) or printed wiring board (PWB) is the object that connects and interfaces most electronic components with each other and with other elements in computers, communication devices, consumer electronics, automated manufacturing and inspection equipment.\nThe procedures of manufacturing these circuit boards and of inserting and connecting multiple components, such as resistors, capacitors and integrated circuits, can be applied in mass production environments, achieving substantial automation, which results in costs reduction, high reliability and high component packaging densities. Backplanes and panels (interconnecting boards, in which printed circuits, panels or integrated circuit packages can be plugged or mounted into or onto) are also manufactured in a similar manner. Modern, highly-dense, populated boards require sophisticated and high resolution manufacturing techniques with precise registration capabilities.\nA PCB is produced from a base of insulating material on which a thin copper layer is laminated or plated, known as a bare copper plated board, from which a chemical etching step selectively removes areas of the copper to produce electrically conducting paths. This selective removal is achieved by covering the copper layer with a patterned mask (etch-resist) that protects the copper layer in the following etching step. For simpler PCBs, screen-printing techniques are generally utilized to form the patterned mask, and for more densely populated PCBs having generally complicated multi-layer conducting paths, Liquid Photo Imageable etch and solder resist mask procedures are commonly utilized. The pattern that remains on the board after the etching step is commonly known as the primary image conductor pattern.\nThe components leads must then be connected to predetermined positions in the conducting paths (called pads) by soldering the leads and the conducting paths utilizing a molten metal alloy, which, after solidifying, achieves a permanent electrically conductive bond. In mass production, wave-soldering methods are commonly utilized, wherein the PCB passes through a molten solder wave that coats the pads and leads and thus forms the required solder joints. A solder resist mask protects the conducting paths from being coated with solder during the soldering step. The solder resist mask leaves uncovered only the pads that need to be covered by the molten solder; otherwise, the conducting paths would also be covered with solder, causing several problems such as short circuits by bridging solder.\nVarious finishing processes in the manufacturing of PCBs currently use (or will be able to use shortly) state of the art inkjet printers. These printers are used to deposit specific kinds of material on the surface of the PCB according to computer generated graphics.\nIn one application, an inkjet printer may be used to cover the PCB with a material that acts as a solder mask in the ensuing manufacturing processes and acts as a protective cover layer throughout the life of the finished electronics subsystem. This solder mask may be applied using any conventional method such as “curtain coating” and “silkscreen” methods.\nThe solder mask coating may also be applied digitally, as described, for example, in U.S. Patent Application Publication No. 2005/0176177 (Zohar et al), the disclosures of which are incorporated herein in their entireties, and which disclose a method and apparatus for applying ink, according to a solder mask pattern, to a printed circuit board having elevated pads defining pad edges. The method includes flooding the printed circuit board with ink such that the ink advances to the pad edges and is stopped thereby and thereat, without climbing onto the elevated pads.\nIn another application, inkjet printers deposit notation ink to print human or machine readable information such as legend, 1D or 2D Barcode, geometry defining lines, etc. See, for example, U.S. Pat. No. 6,754,551 (Zohar et al.), the disclosure of which is incorporated herein in its entirety, and which describes a jet dispensing print system for dispensing a liquid or viscous substance as a pattern onto the surface of a PCB in an industrial manufacturing PCB production line.\nThere is a growing need to speed up the manufacturing process of PCBs, to reduce the footprint of PCB manufacturing tools and to reduce the cost associated with the manufacturing of PCSs."}
{"text": "This application is based on Japanese Patent Application Nos. 11-153060 (1999) filed May 31, 1999, 11-153062 (1999) filed May 31, 1999, 11-153063 (1999) filed May 31, 1999, 11-153064 (1999) filed May 31, 1999, and 2000-117063 filed Apr. 18, 2000, the contents of which are incorporated hereinto by reference.\n1. Field of the Invention\nThe present invention relates to an ink tank, an ink-jet cartridge, an ink-supplying apparatus, an ink-jet printing apparatus, a method for supplying ink, an ink-jet printing head and a printing apparatus.\n2. Description of the Related Art\n(First Prior Art)\nHeretofore, a serial-scanning type printing apparatus has been known as an example of the ink-jet printing apparatus. This kind of the printing apparatus exchangeably carries a printing head as a printing means and an ink tank as an ink container on the carriage which is capable of movement in the direction of main-scanning perpendicular to the direction of sub-scanning (i.e., the direction of moving a printing medium such as a piece of paper). As for this kind of the printing system, images are sequentially printed on a printing medium by repeating the movement of the carriage on which the printing head and the ink tank are mounted in the direction of main-scanning and the movement of the printing medium in the direction of sub-scanning.\nThe serial-scanning type printing apparatus is able to print an image on a large sized printing medium (e.g., A1, A0 size) by enlarging the migration width of the carriage. In this case, however, the ink storage capacity of the ink tank should be increased for using a great volume of ink to print an image on the surface of a large-sized printing, so that the whole weight of the carriage is increased in proportion to the capacity of the ink. In addition, an inertial force in the movement of the carriage is also proportionally increased. For moving the carriage at a high speed against the inertial force, there is the need for installing a driving motor with a large amount of electric power for driving the carriage in high power, resulting in the problem of increasing the price of the printing apparatus in its entirety. In addition, as the total weight of the carriage is increased, there is another problem that the printing apparatus oscillates greatly as a whole by the counterforce contrary to the force for deaccelerating the carriage to zero against the inertial force when the carriage returns at a returning point of its reciprocating motion in the main-scanning direction. Therefore, it was difficult for speeding up the travel speed of the carriage.\nFor reducing the weight of the carriage, on the other hand, the capacity of the ink tank may be lessened. In this case, however, the frequency of replacing the ink tank rises and thus there is a high possibility of replacing the ink tank with the new one in the middle of the printing movement.\nOne of the solutions to solve the problem about such a replacement of the ink tank is proposed in Japanese Patent Application Laying-open 9-24698 (1997). In this prior art document, a deformable ink container is connected to a printing head. The deformable ink container can be connected to an auxiliary ink container as necessary for supplying ink from the latter to the former. The deformable ink container comprises a bag that stores ink under the negative pressure enough to restrain the leakage of ink from the ink-eject port. Therefore, ink can be supplied from the auxiliary ink container to the deformable ink container by an effect of such a negative pressure.\nThe bag used in the deformable ink container is a flexible one enough to reduce its capacity in proportion to become flat, depending on the volume of ink ejected from the printing head (i.e., the usage of ink in the bag). When the volume of the bag is decreased to less than the fixed volume, a supply opening of the deformable ink container is opened to establish connection with the auxiliary ink container. As a result, ink is supplied into the bag of the deformable ink container from the auxiliary ink container by the negative pressure of the inside of the bag. When the ink capacity of the bag reaches to a maximum level, the negative pressure in the bag becomes zero and the supply of the ink is automatically stopped. According to such a prior art, therefore, the supply of ink can be automatically stopped by using the negative pressure without requiring the control using a pressure sensor, a volume detection sensor, and so on.\nBy the way, the upper limit of the negative pressure in the deformable ink container can be determined by its balance with the force of ejecting ink from the printing head. If the negative pressure becomes too high, the force of ejecting ink from the printing head is decreased by an effect of the negative pressure. Therefore, the negative pressure must be decided within the scope of the best ink-eject conditions in the printing head. In addition, a head location of ink in the auxiliary ink container must be configured so that it is lower than that of ink in the deformable ink container. If the deference between those heads is too large, ink cannot be supplied any more even if the negative pressure in the deformable ink container is defined so as to correspond to the conditions of ink-eject of the printing head.\nAs for the prior art, therefore, it is provided with the special device to configure a position of the auxiliary ink container in the vertical direction with respect to the deformable ink container. As for being provided with such a device, however, the problems of upsizing and cost up of the printing apparatus may be caused. If air enters into an ink flow path that connects between the auxiliary ink container and the deformable ink container from a part of the path at the time of ink supply, the entering air moves into the bag of the deformable ink container and then reduces the ink capacity of the deformable ink container by a large amount. Furthermore, the deformable ink container is filled with air if a large amount of the air is entered into the bag, so that there is a problem that a further supply of ink cannot be made. Still furthermore, the deformable ink container comprises an elastic container part that forms a bag and a movable part such as a spring that inflate the bag to a predetermined volume. Thus, there are further problems of the limitation of downsizing, complicated and heavy-weighted structure, and the rise in production cost.\n(Second Prior Art)\nHeretofore, a serial-scanning type printing apparatus has been known as an example of the ink-jet printing apparatus. This kind of the printing apparatus exchangeably carries a printing head as a printing means and an ink tank as an ink container on the carriage which is capable of movement in the direction of main-scanning perpendicular to the direction of sub-scanning (i.e., the direction of moving a printing medium such as a piece of paper). The printing head and the ink tank are connected each other by an ink path. As for this kind of the printing system, images are sequentially printed on a printing medium by repeating the movement of the carriage on which the printing head and the ink tank are mounted in the direction of main-scanning and the movement of the printing medium in the direction of sub-scanning.\nOn the other hand, a method for supplying ink to the ink tank of the ink-jet printing apparatus may be of the supply of ink through the application of pressure to the ink or the sucking of ink through the induction of negative pressure in the ink tank.\nBy the way, if the method for sucking of the ink into the ink tank is used as a method for supplying ink to the ink tank being connected to the printing head, there is the possibility of sucking ink in the printing head into the ink tank by an effect of the negative pressure to be introduced into the ink tank at the time of supplying ink under suction. If the ink in the printing head is introduced into the ink tank, a meniscus of ink to be formed on each of ink eject ports of the printing head is broken down and air enters into the printing head through the ink eject port. As a result, the supply of ink under suction cannot be performed as the negative pressure in the ink tank is reduced.\n(Third Prior Art)\nHeretofore, a printing apparatus that performs the printing using a printing material such as ink have been widely available. In recent years, in particular, a serial-scan type ink-jet printing apparatus is rapidly becoming in widespread use. Such an ink-jet printing apparatus comprises a carriage on which a printing head and an ink tank are mounted. The printing head ejects ink onto a printing medium to print an image thereon while the carriage moves directly above the printing medium in the main-scanning direction.\nAccording to the configuration of such a printing apparatus, an empty ink tank must be replaced with the new one to continue its printing movement when the ink stored in the ink tank is exhausted. If the printing movement is continued long or performed on a larger-sized printing medium, a larger amount of ink may be consumed. In this case, therefore, the ink tank must be exchanged frequently, so that the printing movement in progress is suspended every time the ink tank is replaced with the new one. Such a replacement work is very troublesome.\nAs a consequence, there is another printing apparatus having a supplementary ink tank for automatically refilling ink when the ink tank mounted on the carriage becomes empty. The supplementary ink tank is connected to the ink tank on the carriage through a tube or the like. Ink can be supplemented from the supplementary ink tank to the ink tank on the carriage when the amount of ink stored in the ink tank decreases to a predetermined level. Therefore, the user may only replace the supplementary ink tank with the new one.\nThe conventional supplementary ink tank generally comprises an ink bag for storing ink and a case for encasing the ink bag.\nThe ink bag may be formed as the joining of two thin films by welding their opposite sides together or by any of other conventional techniques. Each of the thin films is generally in the shape of a rectangular, and also a part of one joining side of the rectangular is shaped like a cylinder as a protrusion being connected to a cylindrically shaped withdrawal member made of plastic or the like. Thus, the ink bag can be fixed in the inside of the case by putting the withdrawal member into an ink output opening of the case.\nA main body of the printing apparatus has a hollow tube that has an external diameter enough to be inserted into the withdrawal member. If the supplementary ink tank is inserted into the predetermined position in the printing apparatus, the hollow tube fits into the withdrawal member of the ink bag and then the connection between the supplementary ink tank and the hollow tube is accomplished. Consequently, the ink tank on the carriage is able to receive ink passing through the hollow tube.\nAlternatively, the supplementary ink tank may be prepared by welding thin films so that the ink bag itself has a cylindrical protrusion without installing any withdrawal means on the ink bag. In this case, the insertion of a needle-like tip of the tube into the protrusion of the ink bag allows the connection between the protrusion and the hollow tube for forming an ink passage.\nHowever, the above conventional supplementary ink tank has the following programs.\nThat is, if a part of the ink bag is formed as a protrusion, the process of shaping the ink bag is complicated and the cost of production is increased.\nFurthermore, if the ink passage between the withdrawal member and the hollow tube is not securely formed, leakage of ink might occur from the loosely connected portion. For automatically connecting them to make an ink passage at the time of mounting the supplementary ink tank, the supplementary ink tank must be precisely connected to the hollow tube so that a center of the withdrawal member coincides with an extension line of a center of the hollow tube. In this case, however, it is difficult to keep such an ink-passage connection consistently because there is a possibility that the hollow tube is curved by putting in and out the supplementary ink tank over and over again.\nIt is a first object of the present invention is to provide an ink tank, an ink-jet cartridge, an ink-supplying apparatus, an ink-jet printing apparatus, and a method for supplying ink, where ink can be reliably supplied to the ink tank by a simplified configuration of an ink passage to achieve both the size and weight reductions of the printing apparatus and to increase the reliability thereof.\nIt is a second object of the present invention is to provide an ink tank, an ink-jet cartridge, an ink-supplying apparatus, an ink-jet printing apparatus, and a method for supplying ink, where ink can be smoothly supplied during an extended period of time.\nIt is a third object of the present invention is to provide an ink-jet printing apparatus, an ink-supplying apparatus, and a method for supplying ink, where ink can be reliably supplied to the ink tank by a simplified configuration of an ink passage to achieve both the size and weight reductions of the printing apparatus and to increase the reliability thereof.\nIt is a fourth object of the present invention is to provide an ink tank, an ink-jet printing head, an ink-jet cartridge, and an ink-jet printing apparatus, where ink can be reliably supplied to the ink tank by preventing the entry of ink or air from the ink-jet printing head connected to the ink tank when ink is supplied to the ink tank under suction caused by the induction of negative pressure in the ink tank.\nIt is a fifth objet of the present invention is to provide an ink tank and a printing apparatus, where the ink tank has a main body that can be easily shaped like a bag and connected to an ink passage at the time of mounting the ink tank on the printing apparatus.\nIn the first aspect of the present invention, there is provided an ink tank capable of introducing ink into the ink tank through an inlet by a negative pressure introduced into the ink tank through a suction port, comprising gas-liquid separating means which is provided at the suction port and which permits gas to pass but inhibits ink from passing.\nIn the second aspect of the present invention, there is provided an ink-jet cartridge comprising an ink tank according to the first aspect, and an ink-jet printing head which is able to eject ink introduced from the ink tank.\nIn the third aspect of the present invention, there is provided an ink-supplying device for supplying ink to an ink tank according to the first aspect or an ink tank of an ink- jet cartridge according to the second aspect, comprising ink-supplying means for supplying ink stored in a main ink tank into the ink tank through the inlet, and negative-pressure loading means for loading negative pressure caused by a suction pump into the ink tank through the suction port.\nIn the fourth aspect of the present invention, there is provided an ink-supplying device for supplying ink to an ink tank according to the first aspect or an ink tank of an ink-jet cartridge according to the second aspect, comprising ink-supplying means for supplying ink stored in a main ink tank into the ink tank through the inlet, negative-pressure loading means for loading negative pressure caused by a suction pump into the ink tank through the suction port, and capping means capable of capping an ink eject port of the printing head by a cap member.\nIn the fifth aspect of the present invention, there is provided an ink-jet printing apparatus, comprising a mounting portion on which an ink tank according to the first aspect and an ink-jet printing head are mountable, where the ink-jet printing is able to eject ink supplied from the ink tank, and transfer means which performs the relative movements of the ink-jet printing head and a printing medium.\nIn the sixth aspect of the present invention, there is provided an ink-jet printing apparatus, comprising a mounting portion on which an ink-jet cartridge according to the second aspect is mountable, and transfer means for relatively moving the ink-jet cartridge and a printing medium.\nIn the seventh aspect of the present invention, there is provided a method for supplying ink to an ink tank according to the first aspect and an ink tank of an ink cartridge according to the second aspect, comprising the steps of supplying ink into the ink tank from the inlet by loading negative pressure into the ink tank from the suction port through the gas-liquid separating means, and stopping the load of negative pressure into the ink tank from the suction port.\nIn the eighth aspect of the present invention, there is provided an ink-jet printing apparatus, comprising a mounting portion on which an ink tank according to the first aspect and an ink-jet printing head are mountable, where the ink-jet printing is able to eject ink supplied from the ink tank, transfer means which performs the relative movements of the ink-jet printing head and a printing medium, and means for forming ink meniscus on the ink eject port by the recovery process which discharges ink from the ink eject port of the ink-jet printing head under suction before supplying of ink to the ink tank.\nIn the ninth aspect of the present invention, there is provided an ink-jet printing apparatus for printing an image on a printing medium employing an ink-jet printing head capable of ejecting ink supplied from an ink tank, comprising negative-pressure loading means which is able to introduce negative pressure into the ink tank, ink-supplying means for supplying ink into the ink tank using the negative pressure in the ink tank, gas-liquid separating means which lies in a negative-pressure loading passage between the ink tank and the negative-pressure loading means and which permits gas to pass but inhibits ink from passing, and disrupting means capable of disrupting a midcourse portion of the negative-pressure loading passage between the ink tank and the gas-liquid separating means.\nIn the tenth aspect of the present invention, there is provided an ink-supplying device, comprising negative-pressure loading means which is able to introduce negative pressure into an ink tank, ink-supplying means for supplying ink into the ink tank using the negative pressure in the ink tank, gas-liquid separating means which lies in a negative-pressure loading passage between the ink tank and the negative-pressure loading means and which permits gas to pass but inhibits ink from passing, and disrupting means capable of disrupting a midcourse portion of the negative-pressure loading passage between the ink tank and the gas-liquid separating means.\nIn the eleventh aspect of the present invention, there is provided a method for supplying ink to an ink tank, comprising gas-liquid separating means which lies in a negative-pressure loading passage between the ink tank and the negative-pressure loading means and which permits gas to pass but inhibits ink from passing, and disrupting means for disrupting a midcourse portion of the negative-pressure loading passage between the ink tank and the gas-liquid separating means, the method comprising the steps of loading negative pressure into the ink tank through the negative-pressure loading passage, supplying ink into the ink tank using negative pressure in the ink tank, stopping the loading of negative pressure into the ink tank by the gas-liquid separating means when ink touches the gas-liquid separating means, and disrupting the midcourse portion by the disrupting means except when ink is supplied into the ink tank.\nIn the twelfth aspect of the present invention, there is provided an ink tank which has an ink-supplying port for supplying ink into an ink-jet printing head, and which is capable of introducing ink into the ink tank by negative pressure introduced into the ink tank, comprising a valve provided at the ink-supplying port, which closes the ink-supplying port by negative pressure higher than a predetermined level in the ink tank.\nIn the thirteenth aspect of the present invention, there is provided an ink-jet printing head capable of ejecting ink supplied from an ink tank through an ink supplying port, comprising a valve provided at a connecting port connected to the ink-supplying port, which closes the ink-supplying port by negative pressure higher than a predetermined level in the ink tank.\nIn the fourteenth aspect of the present invention, there is provided an ink-jet cartridge comprising an ink tank according to the twelfth aspect, and an ink-jet printing head capable of ejecting ink supplied from an ink tank through an ink-supplying port.\nIn the fifteenth aspect of the present invention, there is provided an ink-jet cartridge comprising an ink-jet printing head according to the thirteenth aspect, and an ink tank capable of supplying ink into the inkjet printing head through the connecting port.\nIn the sixteenth aspect of the present invention, there is provided an ink-jet printing apparatus comprising a tank mounting portion on which an ink tank according to the twelfth aspect is mountable, a head mounting portion on which an ink-jet printing head capable of ejecting ink supplied from the ink tank is mountable, and moving means for relatively moving the ink-jet printing head and a printing medium.\nIn the seventeenth aspect of the present invention, there is provided an ink-jet printing apparatus comprising a head mounting portion on which an ink-jet printing head according to the thirteenth aspect is mountable, a tank mounting portion on which an ink tank capable of supplying ink to the ink-jet printing head is mountable, and moving means for relatively moving the ink-jet printing head and a printing medium.\nIn the eighteenth aspect of the present invention, there is provided an ink tank having a bag-like tank body which is made of a sheet of a thin film that is folded down in one side to form a folding part, and which is capable of storing ink, wherein the folding part forms a connecting portion capable of connecting between the inside and the outside of the tank body by means of a hollow conduit that is able to penetrate the folding part.\nIn the nineteenth aspect of the present invention, there is provided a printing apparatus capable of printing of an image using ink in the tank body, comprising a tank mounting portion on which an ink tank according to the eighteenth aspect is mountable, wherein a hollow conduit that is able to penetrate the connecting portion of the tank body and is provided at the tank mounting portion.\nThe present invention is configured such that the supply of ink under suction can be automatically stopped using the function of a gas-permeable member, so that the supply of ink to the ink tank can be performed by a simple structure with reliability. This offers an advantage of being able to achieve both the size and weight reductions of the printing apparatus and an improved reliability thereof.\nThe present invention is also configured such that the formation of ink meniscus on an ink eject port of the printing head is performed by draining the ink from the printing head being connected to the ink tank under suction, before the supply of ink to the ink tank under suction is performed. This offers an advantage of being able to achieve the supply of ink to the ink tank under suction with reliability.\nThe present invention is configured such that a porous material with an oil repellent finish is used as the gas-permeable member to be functioned as a gas-liquid separate means. The gas-permeable member repels ink enough. This offers an advantage of being able to achieve the supply of ink smoothly over an extended period of time with reliability in addition to improve the durability of the gas-permeable member.\nThe present invention is configured such that the gas-liquid separate means is not connected to the inside of the ink tank except when the supply of ink is performed. This offers an advantage of being able to prevent that the performance of the gas-liquid separate means is decreased by exposing the gas-liquid separate means to ink for a long time.\nThe present invention is configured such that a valve is provided in an ink-supplying path between the ink tank and the ink-jet printing head and closed when the inside of the ink tank becomes a predetermined level of negative pressure. This offers an advantage of being able to achieve the supply of ink under suction with reliability by preventing the entry of ink or air from the ink-jet printing head to be connected to the ink tank.\nThe present invention is configured that the main body of the ink tank is shaped like a bag which is good enough for communicating the inside of a main body of the ink tank with the outside through a hollow tube by passing the hollow tube through a curved portion of a thin film that forms the bag-shaped main body of the ink tank. This offers an advantage of being able to achieve the cost reduction of manufacturing the ink tank as the bag-shaped main body of the ink tank is formed with ease.\nThe above and other objects, effects, features and advantages of the present invention will become more apparent from the following description of embodiments thereof taken in conjunction with the accompanying drawings."}
{"text": "The present invention relates to power supply apparatus and methods and, more particularly, to apparatus and methods for transferring loads among multiple power sources.\nA “reverse transfer” or “on-line” uninterruptible power supply (UPS) may be used to provide protected power to a load. As shown in FIG. 1, a typical on-line UPS 110 includes a rectifier 112 that is configured to be coupled to an AC power source 10 and to produce a DC voltage therefrom. The DC voltage is applied to an inverter 114, which generates an AC voltage for a load 20 coupled thereto. An alternative power source, e.g., a battery 116, may be coupled to the intermediate DC link to provide power to the inverter 114 in the event of failure of the AC power source 10. If the UPS 110 fails or is taken off line for maintenance, a high-speed solid-state static switch 120, which, as shown, may include anti-parallel connected silicon controlled rectifiers (SCRs) or other solid-state circuits that provide similar switching capabilities, may be used to provide a bypass path between the power source 10 and power to the load 20. The high-speed nature of the static switch 120 allows the load 20 to be transferred to the main source 10 with little or no interruption. In some applications, the static switch 120 may be integrated with the UPS 110 while, in other applications, the static switch 120 may be a separate device. As shown, a lower-speed switching device 130, such as switch, relay or circuit breaker, may be provided to allow for bypassing of the static switch 120. Other switches/breakers 118 may be provided for isolation and/or circuit protection.\nStatic switches may also be used to provide for transfer of loads among UPSs. Referring to FIG. 2, a first UPS 110a, here shown as including bypass static switch and other isolation/circuit protection devices, may be coupled to a first power source 10a and switchably coupled to a first load 20a via a static switch 212 of a first static transfer switch (STS) 210a and to a second load 20b via a static switch 212 of a second STS 210b. Similarly, a second UPS 110b may be coupled to a second power source 10b and switchably coupled to the first load 20a via a static switch 212 of the first STS 210a and to the second load 20b via a static switch 212 of the second STS 210b. A switch/breaker 220 may be provided between the outputs of the UPSs 110a, 110b. \nThe STSs 210a, 210b allow for transfer of the loads 20a, 20b between the first and second UPSs 110a, 110b by providing for selective coupling therebetween. In particular, at a given time, one static switch 212 in each STS 210a, 210b is activated to provide power to one of the loads 20a, 20b, while the other static switch 212 isolates the same load from the other power source. Each of the STSs 210a, 210b can transfer a load in an uninterrupted manner by substantially simultaneously turning off the active static switch and turning on the inactive static switch. The STSs 210a, 210b are typically built as integrated units that include electronic circuits to monitor voltages applied to the loads 20a, 20b and responsively control the static switches 212. It will be appreciated that, although FIG. 2 illustrates dual STSs 210a, 210b, STSs with more than two static switches 212 may be used in applications in which more than two UPSs (or other power sources) are to be interconnected. A relatively less complex load transfer capability has also been provided in conventional systems by using a single static transfer switch that is coupled between two UPS outputs (or other power sources) and that operates responsive to a voltage at one or both of the outputs."}
{"text": "Secondary batteries typically comprise an electrode assembly in which a positive electrode plate and a negative electrode plate are spirally wound with a strip-shaped separator interposed therebetween. The positive electrode plate includes a strip-shaped current collector and a positive electrode material mixture layer formed on the current collector. The negative electrode plate includes a strip-shaped current collector and a negative electrode material mixture layer formed on the current collector. The innermost winding turn of the spirally-wound electrode assembly has an exposed portion of the positive electrode current collector on which no positive electrode material-mixture layer is formed. In the structure called a positive electrode lead forward structure, a positive electrode lead is connected to the exposed portion of the positive electrode current collector, in other words, the positive electrode lead is connected to the winding start portion of the electrode assembly.\nIn order to prevent an internal short-circuit from occurring in the electrode assembly having such positive electrode lead forward structure, an arrangement is proposed in which the distance between the positive electrode lead and the negative electrode plate is increased (see, for example, Patent Document 1 and Patent Document 2).    Patent Document 1: Japanese Patent No. 3237015    Patent Document 2: Japanese Patent No. 3373934"}
{"text": "1. Field of the Invention\nThis invention relates to photography, and more particularly, it relates to compositions and processes useful in the development of photosensitive silver halide emulsions.\n2. Description of the Prior Art\nIn recent years, there has been a growing interest in heterocyclic developing agents, i.e., developers containing a heterocyclic ring as part of their structure. Some of these developing agents have the conventional hydroxyl or amino developing groups substituted on adjacent carbon atoms of a heterocyclic ring to provide structures similar to those of the developing agents in the aliphatic and aromatic series. Still other heterocylic developing agents have one of the functional developing groups included as part of the heterocyclic ring. Illustrative developing agents of the latter type include 1-phenyl-3-pyrazolidone and its 4,4-dialkyl derivatives which form the subject matter of U.S. Pat. No. 2,289,367 and 2,772,282, respectively, and 1-phenyl-4-amino-5-pyrazolones which form the subject matter of U.S. Pat. No. 2,895,825.\nThe present invention is concerned with the use of certain 4-hydroxy pyrazoles as silver halide developing agents."}
{"text": "1. Field of the Invention\nThis invention relates to an improved remote control system for a combined ceiling fan and light fixture and, more particularly, to a system of this type using a radio frequency link, microprocessor control, and a number of special features.\n2. Description of the Prior Art\nCeiling fans are often combined with light fixtures, and it is desirable to provide remote control means for controlling fan speed and direction as well as light intensity. One such remote control system is shown in U.S. Pat. No. 4,719,446, issued Jan. 12, 1988 to Edward F. Hart. This system uses existing house wiring to couple control signals from a wall mounted control unit to control the speed of the fan and the intensity of the light. However, systems of this character require access to an AC line by replacing a wall switch and are subject to line noise. Such systems are difficult to operate in that fan and light operations are sometimes confused. Moreover, the prior art systems do not readily lend themselves to the incorporation of special features. Fan motors of the prior art are typically connected across an AC power supply through triacs which, in conventional practice, are triggered continuously. The commutating effect of back EMF from the motor causes the triac to go to zero at some point other than true AC zero cross, thus causing the triac to fall out of conduction. Although continuous triggering is effective, the prior art fan motors tend to have high power dissipation causing the motor to run hot."}
{"text": "Ureteral stents are used to assist urinary drainage from the kidney to the bladder in patients with ureteral obstruction or injury, or to protect the integrity of the ureter in a variety of surgical manipulations. More specifically, stents may be used to treat or avoid ureter obstructions (such as ureteral stones or ureteral tumors) which disrupt the flow of urine from the kidneys to the bladder. Serious obstructions may cause urine to back up into the kidneys, threatening renal function. Ureteral stents may also be used after endoscopic inspection of the ureter.\nUreteral stents typically are tubular in shape, terminating in two opposing ends: a kidney (upper) end and a bladder (lower) end. The ends may be coiled in a pigtail or J-shape to prevent the upward or downward migration of the stent, e.g., with physiological movements. The kidney coil is designed to retain the stent within the renal pelvis of the kidney and to prevent stent migration down the ureter. The bladder coil sits in the bladder and is designed to prevent stent migration upwards toward the kidney. The bladder coil is also used to aid in retrieval and removal of the stent.\nUreteral stents, particularly the portion positioned in the ureter near the bladder and inside the bladder, may produce adverse effects including blood in the urine, a continual urge to urinate, strangury, and flank pain accompanying reflux of urine up the stent (e.g., when voiding) as pressure within the bladder is transmitted to the kidney. In short, stents may cause or contribute to significant patient discomfort and serious medical problems.\nFIG. 10 is a schematic drawing of the human urinary tract without a stent, showing the renal pelvis, the kidney, the ureter, and the ureteral orifices opening into the bladder. FIG. 11 depicts a typical double-J stent 10 which comprises a small tube 12 which sits inside the urinary system and assists the flow of urine from the kidney (renal pelvis) to the bladder. FIG. 12 depicts prior art indwelling ureteral stent 10 in position. Such stents are typically made of biocompatible plastic, coated plastic, or silicone material. Tube 12 typically varies in size from 4-8 fr. (mm in circumference), and it has multiple small holes throughout its length. A coiled shape pre-formed at each end 14 and 16 is designed to confine its movement within the urinary system, so that it will be maintained in the desired position. The upper (kidney) end 14 of the stent may be closed or tapered, depending on the method of insertion (e.g., the use of a guidewire). The tubular stent extends through the ureteral orifice 18a and into the bladder, fixing orifice 18a open, and thereby enhancing the opportunity for reflux. For clarity, the ureter entering bladder 20 through orifice 18b is not shown. A monofilament thread 22 may be attached to the bladder end of the stent for removal, usually without cystoendoscopy.\nU.S. Pat. No. 4,531,933 (“the '933 patent”) discloses a ureteral stent having helical coils at each end which are provided for preventing migration and expulsion."}
{"text": "The micro-movement screw subassembly for linear displacement is assembled by screws and screw nuts, and the movement for lead and gear is combined together. However, a gap error always exists between the screw and screw nut, which is caused by error in machine processing for triangular thread surfaces and axial face as well as screw pitch deviation. Generally, when the error is more than 10 μm, it causes inaccurate screw leading linear displacement. Conventionally, the first solution is to use organic powder during the processing, and the second solution is to open an elastic groove on the screw shaft to reduce the gap error between the screw and the screw nut. The first solution is not easy and convenient, as the powder used to fill in the gap continuously falls out and needs to be replenished during work. The second solution does not work well either, while gap is reduced but the leading accuracy is still not good."}
{"text": "The present invention relates to a beverage preparation machine, and particularly though not exclusively to an espresso-type coffee machine that includes an outlet nozzle for steam or other liquids such as milk. These nozzles, called steam nozzles, are immersed into the beverage. They subsequently require a cleaning of at least their outside surface that has come in contact with the beverage. More particularly, the invention relates to a beverage preparation machine comprising:                a body;        a movable steam nozzle connected to the body by a moving device;        a beverage preparation area in which a cup collecting the beverage can be placed;        a cleaning area which is separate from the preparation area and in which is disposed at least one rinsing container for the nozzle, the moving device being capable of moving the nozzle between the preparation and cleaning areas.        \nCleaning a steam nozzle, for example after preparing a cappuccino or a latte, is time-consuming since conventionally, it requires a wet cloth or a container of water and cleaning product. Consequently, after the preparation of a beverage, the steam nozzle is quite often left covered with residue, particularly in the case of a home machine used by several people. Cleaning it becomes even more difficult after traces of milk have dried on the outside wall of the nozzle. It is then practically essential to let the nozzle soak in a container containing a cleaning solution, typically water mixed with a cleaning product like a detergent. Yet the presence of a container, a source of hot water if the machine does not have one, and particularly a bottle of cleaning product near the machine is not very satisfactory. Aside from the space issues, it is necessary to make sure that the bottle of cleaning product is actually present, that it still contains some liquid, that this liquid is measured out correctly by the users, and that it is not within reach of children. To avoid altering the flavor or appearance of a beverage prepared afterwards, it is also important to thoroughly rinse or wipe off the nozzle after it is cleaned.\nAll these inconveniences mean that the cleaning of a steam nozzle is often neglected, even though in certain machines it is relatively easy to manually move the steam nozzle in order to immerse it into a receptacle placed next to the cup rest.\nMoreover, there are systems, known particularly from the document WO03091152A, for automatically cleaning a nozzle similar to a steam nozzle. These systems place the nozzle in a relatively closed chamber and dispense a cleaning solution under pressure through the nozzle. The vigorous agitation of the liquid around the nozzle cleans it. However, the main purpose of these systems is to clean out the internal conduits through which a product, like milk, that is susceptible to bacterial contamination has circulated. These systems require solenoid valves and complex and expensive dosing devices. Aside from having to switch the nozzle between a beverage supply and a supply of water under pressure, it is necessary to provide for the delivery into the conduits of a cleaning product in precise doses, on the order of a few percent of the water flow. These systems operating under pressure are therefore costly, especially since it is necessary to guarantee high degrees of watertightness and operational reliability in order to avoid dispensing a residue of cleaning solution along with the beverage."}
{"text": "1. Field of the Invention\nThe present invention relates to compounds which inhibit botulinum neurotoxin. In particular, the present invention relates to compounds which inhibit botulinum neurotoxin serotype A light chain (BoNT/A LC) metalloprotease activity.\n2. Description of the Related Art\nBotulinum neurotoxins (BoNTs) are produced by spore forming anaerobic bacteria Clostridium botulinum, and are among the most lethal of biological poisons (LD50=0.001 μg per Kg). See Schmidt & Stafford (2003) Appl. Environ. Microbiol. 69:297-303; Kessler & Benecke (1997) Neurotoxicology 18:761-770; and Burnett et al. (2005) Nat Rev Drug Discov 4(4):281-297. Seven immunologically distinct BoNT serotypes (designated A-G) have been identified. See Simpson, L. L. (1989) BOTULINUM NEUROTOXIN AND TETANUS TOXIN, Academic Press, New York.\nExposure to BoNTs, for example, through contaminated food, can result in life threatening flaccid paralysis. See Shapiro, et al. (1998) Ann. Intern. Med. 129:221-228. Furthermore, BoNTs have been weaponized in highly toxic aerosol form, and consequently pose a significant threat to both to civilian and military populations. See Franz, et al. (1997) JAMA 278:399-411; and Amon, et al. (2001) JAMA 285:1059-1070.\nAs indicated, these enzymes have been weaponized in aerosol medium, and airborne release or direct contamination (e.g. foodstuffs) represent significant threats to both military and civilian populations. See Paddle, B M (2003) J Appl Toxicol 23(3):139-170; Clarke, S C (2005) Br J Biomed Sci 62(1):40-46; Hicks et al. (2005) Curr Med Chem 12(6):667-690; and Josko, D (2004) Clin Lab Sci 17(1):30-34. And, with the increased use of BoNTs as therapies for a range of medical conditions and superficial cosmetic purposes, there is the increased potential for accidental overdosing. Furthermore, as the popularity of BoNTs as therapeutics continues to grow, these enzymes are increasingly being manufactures overseas, where less strict controls may allow clandestine organizations to obtain large quantities of these toxins in very concentrated, pure, and easily stored formulations. See Comella & Pullman (2004) Muscle Nerve 29(5):628-644; Gormley et al. (1997) Muscle Nerve Suppl 6:S14-20; Marks, J D (2004) Anesthesiol Clin North America 22(3):509-532, vii; Noonan & Adler (2002) Newsweek 139(19):50-56, 58; O'Brien, C F (2001) Adv Neurol 87:265-269; O'Brien, C F (2002) Clin J Pain 18(6 Suppl):S182-190; O'Brien, D (2003) J Perianesth Nurs 18(2):126-134; Rossetto (2001) Toxicon 39(1):27-41; Shukla & Sharma (2005) Crit Rev Microbiol 31(1):11-18; and Turton et al. (2002) Trends Biochem Sci 27(11):552-558. As a result, there is an urgent need for therapeutic countermeasures against BoNTs. See Goodnough, et al. (2002) FEBS Lett. 513:163-168.\nBoNT is secreted as a holotoxin composed of two peptide chains that are linked by a disulfide bridge. See Lacy & Stevens (1999) J. Mol. Biol. 291:1091-1104. The heavy chain is responsible for: (1) targeting and binding to surface receptors on nerve terminals; (2) translocation into the neuronal cytosol via the formation of a low pH endosome; and (3) protecting the substrate binding cleft of the light chain prior to neuronal internalization. See Turton, et al. (2002) Trends Biochem. Sci. 27:552-558; and Singh, B. R. (2000) Nat. Struct. Biol. 7 (2000) 617-619. The light chain, which dissociates from the heavy chain in the low endosomal pH, is released into the cytosol where it acts as a zinc metalloprotease that cleaves soluble NSF-attachment protein receptor (SNARE) proteins: synaptosomal-associated protein of 25 kDa (SNAP-25), synaptobrevin, and syntaxin. BoNT serotypes A, C, and E cleave SNAP-25; serotypes B, D, F, and G cleave synaptobrevin; and serotype C can also use syntaxin as substrate. See Binz, et al. (1994) J. Biol. Chem. 269:1617-1620; Schiavo, et al. (1992) Nature 359:832-835; Schiavo, et al. (1993a) J. Biol. Chem. 268:23784-23787; Schiavo, et al. (1993c) J. Biol. Chem. 268:11516-1151915; Schiavo, et al. (1993b) J. Biol. Chem. 269:20213-20216; and Blasi, et al. (1993b) EMBO J. 12:4821-4828. Without functional SNARE complexes, acetylcholine is not released into neuromuscular junctions, thereby leading to paralysis.\nResearch to identify peptide and small molecule inhibitors of BoNT serotype A (BoNT/A) has targeted both holotoxin translocation and light chain (BoNT/A LC) metalloprotease activity. Sheridan et al. and Deshpande et al. have shown that a number of antimalarial agents interfere with BoNT/A translocation into nerve cytoplasm. See Sheridan, et al. (1997) Toxicon 35:1439-1451; and Deshpande, et al. (1997) Toxicon 35:433-445.\nSpecifically, it has been shown that several antimalarial compounds act subsequent to toxin binding to cell-surface receptors, and it has been hypothesized that these agents inhibit BoNT/A cytosol entry by raising endosomal pH (an endosomal pH of 5.5 or lower is needed for release into the cytoplasm). Hayden et al. have found that BoNT/A LC is inhibited by mM concentrations of known protease inhibitors: captopril, lysinopril, and enalapril. See Hayden, et al. (2003) J. Appl. Toxicol. 23:1-7. In the same study, it was also reported that a number of short peptides, from specific “hinge” libraries, inhibit BoNT/A LC activity by as much as 51% at concentrations as low as 0.5 μM. Using a chromatographic method, Schmidt et al. identified the peptide motif CRATKML as a potent inhibitor. See Schmidt, et al. (1998) FEBS Lett. 435:61-64. In a subsequent study, the Cys residue of CRATKML was replaced with thiol containing organic moieties, and it was found that a 2-mercapto-3-phenylpropionly containing derivative was the most effective (Ki=0.3 μM). See Schmidt & Stafford (2002) FEBS Lett. 532:423-426.\nNeither the currently available BoNT antitoxin nor antibodies can counter these toxins once they are inside neurons; currently, critical care mechanical ventilation is the only treatment option. However, the effects of internalized BoNTs can last for months, and mechanical ventilation would be impractical if even a limited number of individuals were simultaneously intoxinated. See Meunier et al. (2003) Mol Cell Neurosci 22(4):454-466; and Eleopra et al. (1998) Neurosci Lett 256(3): 135-138. Furthermore, antitoxin administration would preclude vaccinated individuals from any form of highly beneficial BoNT medical therapy.\nThus, a need exists for small molecule (non-peptidic) inhibitors of BoNT/A LC metalloprotease activity."}
{"text": "THIS invention relates to a corner device to be attached to an article, such as a document formed by a sheet of paper or sheets of paper. This invention further relates to kit for attaching a corner device to a document and a corner device locator to align and hold the corner device before attaching it to the sheets of paper.\nStaples are one of the most commonly used fasteners to secure sheets of paper together. As it is fairly easy to remove a staple which secures sheets of paper together, a need has arisen for a tamper evident device or seal to be applied with the staple, thereby to provide some security features to a document formed by the sheets of paper.\nThe use of a staple to secure sheets of paper together has the disadvantage that the corners of the sheets of paper typically discolour and curl upwards after the document has been used by a number of people or for an extended period of time. This may result in the document, especially in cases where it is a legal document, such as a contract, losing its aesthetic appeal and professional look.\nIt is an object of the present invention to provide a corner device that addresses at least some of these problems."}
{"text": "The present invention relates to compositions for use at start-up a reformer of a fuel cell system. In particular, this invention includes emulsion compositions comprising hydrocarbon fuel, water and surfactant for use at start-up of a reformer of a fuel cell system.\nFuel cell systems employing a partial oxidation, steam reformer or autothermal reformer or combinations thereof to generate hydrogen from a hydrocarbon need to have water present at all times to serve as a reactant for reforming, water-gas shift, and fuel cell stack humidification. Since water is one product of a fuel cell stack, during normal warmed-up operation, water generated from the fuel cell stack may be recycled to the reformer. For start-up of the reformer it is preferable that liquid water be well mixed with the hydrocarbon fuel and fed to the reformer as an emulsion. The current invention provides emulsion compositions suitable for use at start-up of a reformer of a fuel cell system.\nOne embodiment of the invention provides emulsion compositions suitable for use at start-up of a reformer of a fuel cell system comprising hydrocarbon, water and surfactant.\nIn a preferred embodiment, the emulsion composition is a bicontinuous emulsion comprising a coexisting mixture of at least 80 vol % of a water-in-hydrocarbon macro-emulsion and from 1 to 20 vol % of a hydrocarbon-in-water micro-emulsion.\nIn another embodiment of the invention is provided a method to prepare a bicontinuous emulsion comprising a coexisting mixture of at least 80 vol % of a water-in-hydrocarbon macro-emulsion and from 1 to 20 vol % of a hydrocarbon-in-water micro-emulsion comprising mixing hydrocarbon, water and surfactant at low shear.\nIn yet another embodiment is a bicontinuous emulsion composition comprising a coexisting mixture of at least 80 vol % of a water-in-hydrocarbon macro-emulsion and from 1 to 20 vol % of a hydrocarbon-in-water micro-emulsion."}
{"text": "1. Field of the Invention\nThe present invention relates generally to display control apparatus and program, and more particularly to such display and control apparatus and program for displaying both video and character information on a television receiver.\n2. Description of the Related Art\nIn television broadcasts, the practical use of ground-wave television broadcasts as well as BS (Broadcast Satellite) digital broadcasts has been advanced in a digital system. In such digital television broadcast, character information including news flashes and/or subtitles combined in a multiplexed manner is received in addition to the video information. The receiver processes images on the character information as well as the received video information, thereby customizing the display colors, sizes and/or fonts of the characters included.\nIn the past, many proposals have been made for receivers that receive the BS digital broadcasts and ground-wave television broadcasts. For example, Published Unexamined Japanese Patent Application Hei 11-8807 discloses a digital broadcast receiver by which the viewer can easily recognize the content of a program being broadcasted in each channel at present and display in a short time a list of channels that inform the viewer of channel information available. According to this arrangement, when character information representing respective channel programs and titles of all or part of receivable channels is displayed in the list of channels, this list and images of programs being broadcasted at present in certain ones of the channels of the list are displayed.\nPublished Unexamined Japanese Patent Application 2000-50220, discloses a teletext receiver by which the viewer can easily recognize what teletext programs are being broadcasted at present and easily select a desired program. This publication discloses means for displaying a plurality of teletext programs decoded or a plurality of pages of the same program each in a compressed state on a smaller picture than the original display screen and superimposed such that they are shifted sequentially by at least one line.\nPublished Unexamined Japanese Patent Application 2002-199354 discloses a digital broadcast receiver in which the size of a display area for character information and the display position of the character information are changeable. In this arrangement, even when a television broadcast program is displayed in a compressed manner, corresponding credit titles can be displayed in an easy-to-see manner. This arrangement comprises a controller that changes a value of control information annexed to the received character information based on a display position of a video picture of the television broadcast program received, thereby moving the display position of the character information on the picture.\nPublished Unexamined Japanese Patent Application 2002-353829 discloses a ground-wave digital broadcasting receiver and an information acquiring and supplying apparatus in which even when the viewer moves from his or her reception area during viewing a program, he or she is able to continue to view the program that he or she has viewed so far. According to this arrangement, when a ground-wave digital broadcast is received and demodulated, thereby providing its video and voice signals, the viewer acquires station select information from an information source, for example an FM teletext, different from the ground-wave digital broadcast and then selects a desired station based on the acquired station select information.\nPublished Unexamined Japanese Patent Application 2002-353915 discloses a mobile information terminal that performs a schedule management function, a calendar function, a simple-document creating function, an electronic mail function and a function to receive a partial service of the ground-wave digital broadcast. This terminal comprises a tuner that selects a ground-wave digital broadcast and produces a corresponding transport stream, a separator that extracts added information from the transport stream, a generator that generates an electronic program guide image based on the added information, a specifying unit that specifies any one of displayed programs on the electronic program guide picture and a delivering unit that delivers information that specifies the specified program directly or via an externally connected communication device to an output line. Further, it produces a data broadcast image based on the added information.\nPublished Unexamined Japanese Patent Application 2003-161620 discloses a car navigation system that uses a ground-wave digital broadcast to receive area information such as area map information included in the ground-wave digital broadcast. This system comprises a stored-data reader that reads desired map data from map media that have stored map data, a broadcast receiver that receives the ground-wave digital broadcast, a storage device that has stored map data, position sensing means that detects the present position of the user's car using a GPS, a display that displays the map data, an input device by which the user can select information, and a control processor that controls these elements.\nAlthough these prior art technical literatures disclose techniques for displaying video and character information, they fail to refer to image processing for character information in consideration of video information that represents the background of characters to be displayed.\nIn the digital broadcast, credit titles to be displayed in association with the content of the television broadcasting program are arranged to be displayed at a predetermined position on a displayed picture of the television broadcasting program by the teletext including data, characters and credit titles. Especially, in the Published Unexamined Japanese Patent Application 2002-199354 an arrangement is disclosed that solves a problem that when the television broadcasting program and the data broadcast content are displayed in a multi-picture screen, the data broadcast content is difficult to view because it is hidden by the credit titles. The technique automatically calculates a value of a position where the credit titles start to be displayed with the position of a lower left corner of the video picture as a reference in a state in which the video picture of the television broadcasting picture is displayed in a compressed manner, and displays the credit titles near the video picture in an easy-to-see manner. However, this technique does not solve a problem that when displayed, superimposed on the video information that represents the background, the character information will be difficult to view.\nUsually, the character information is displayed, superimposed on the video information that represents the background. Thus, it is difficult to discriminate the character information from its background depending on the content of the background, thereby lowering the distinction of the character information. For example, when the video is similar in display color to the characters or the contrast in brightness between the video and the characters is small, the distinction of characters is low and the characters are difficult to recognize. When a news flash, for example, of a natural disaster or accident is displayed in large characters on the picture to attract the viewers' attention, the picture quality would greatly be deteriorated."}
{"text": "1. Field of the Invention\nThe present invention relates to a pressure sensor apparatus for detecting or measuring a pressure of an object or fluid whose pressure is to be detected or measured such as, for example, a pressure sensor apparatus mounted on an intake manifold of an intake system of an internal combustion engine for detecting the intake air pressure. Hereinafter, the object or fluid whose pressure is to be detected or measured will be referred to as the fluid for measurement only for the convenience of description.\n2. Description of Related Art\nHeretofore, such pressure sensor apparatus is known in which invasion of foreign materials into a pressure sensor constituting a major part of the pressure sensor apparatus is prevented by providing a barrier plate for causing the foreign materials to flow around the pressure sensor. In this connection, reference may have to be made to, for example, Japanese Patent Application Laid-Open Publication No. 30535/1999 (JP-A-1999-30535) (FIG. 1).\nIn the pressure sensor apparatus mentioned above, invasion of the foreign materials into the pressure sensor can certainly be prevented by providing the barrier plate. However, to this end, it is required to prepare a dedicated barrier plate, which gives rise to a problem that the manufacturing cost of the pressure sensor apparatus is correspondingly increased, to a disadvantage."}
{"text": "(a) Technical Field of the Invention\nThe present invention is generally related to heating devices, and more particular to a heating cable control system having an optical coupling circuit, a NTC break-off detection circuit, and a fourth comparator circuit.\n(b) Description of the Prior Art\nAs shown in FIG. 1, U.S. Pat. No. 8,164,035, titled Heating Device Having Dual-core Heating Cable, teaches a dual-core heating cable control system containing a dual-core heating cable 30 and a control device 20. The control device 20 contains a control circuit 6, a DC voltage circuit 22, a first comparator circuit 40, a second comparator circuit 41, a synchronous signal input circuit 23, a first reference voltage circuit 24, an adjustment circuit 27, a NTC (negative temperature coefficient) detection circuit 75, a PTC (positive temperature coefficient) detection circuit 72, a switch circuit, a load detection circuit 28, a protection circuit, a function selection circuit 29, and a status indicator circuit 74. The DC voltage circuit 22 provides DC voltage Vcc to power the control device 20 and, through the switch circuit, to activate the gate of a second silicon-controlled regulator so that the second silicon-controlled regulator conducts its anode and cathode, and that the AC is conducted to a PTC resistive wire of the dual-core heating cable to heat up. The PTC detection circuit obtains a load current through the PTC resistive wire and converts the load current to a voltage compared to the first reference voltage through the second comparator circuit. When a high level is detected, the heating up to the dual-core heating cable continues whereas a low level is detected, the heating up to the dual-core heating cable stops, thereby achieving constant temperature. Additionally, the dual-core heating cable is configured with a NTC resistive layer and electrical current flows through the PTC resistive wire, the NTC resistive layer, and then to the NTC detection circuit. The NTC detection circuit coverts the current to a voltage compared against the first reference voltage through the first comparator circuit. When a high level is detected, the heating up to the PTC resistive wire continues whereas a low level is detected, the heating up to the PTC resistive wire stops, thereby achieving a second over-temperature protection.\nHowever, even though with the constant temperature and the second over-temperature protection, the heating cable control system still suffers the following disadvantage. When the NTC detection circuit stops heating up due to the second silicon-controlled regulator becomes open-circuited. The protection circuit could still trigger a first silicon-controlled regulator to conduct and the PTC resistive wire is still heated up. The NTC detection circuit then cannot accurately detect the breaking off of the NTC resistive layer and top the heating up to the PTC resistive wire. The dual-core heating cable then would be over-heated and damaged, or the user could be burned."}
{"text": "1. Field of the Invention\nThe present invention relates to a composition of polybutene-1 and a process for preparing the same and, more particularly, to a composition and a process for the preparation of polybutene-1 which has less gel content and less fish eyes and which can preferably be used in the field of film.\n2. Description of Related Art\nA homopolymer of polybutene-1 or a copolymer of butene-1 with another olefin, as generally called polybutene-1, has been used for pipes for warm water due to its excellent creep resistance at high temperatures or in the film field due to its high orientation at a low degree of stretching. Such conventional polybutene-1, however, has the drawback that it is remarkably poor in moldability or forming performance because of its low speed of crystallization. Hence, various attempts have so far been made to solve the drawback of conventional polybutene-1.\nIn order to accelerate the cystallization speed, for example, U.S. Pat. No. 4,321,334 proposes using a nucleating agent in preparing polybutene-1. In this process, however, it is so difficult to disperse the nucleating agent in polybutene-1, that it cannot be said to produce the effect to a sufficient degree. Particularly, when polybutene-1 obtained by this process is formed to film, the film suffers from the disadvantage that fish eyes are caused to thereby give a poor appearance.\nFurther, Japanese Patent Unexamined Publication (kokai) No. 123,607/1980 discloses a process for preparing crystalline polybutene-1 having a high bulk density by subjecting a small amount of .alpha.-olefin other than butene-1 to preliminary polymerization in the process for preparing polybutene-1. As this process is developed with the object to provide polybutene-1 with high bulk density by using a small amount of the other .alpha.-olefin for preliminary polymerization, the above-identified patent publication is silent about the problems of gel content and fish eye. Turning now to specific examples of the description of the patent publication, it can be noted that they do not use a molecular weight modifier, such as hydrogen or the like, during preliminary polymerization. Hence, it can be presumed that a polymer obtainable by this preliminary polymerization may have a too high molecular weight, thereby causing gel and fish eye in large numbers."}
{"text": "Mobile cellular networks are typically engineered by operators to provide diverse sets of voice and data services to groups of users sharing network resources. In a standard configuration, each of the networked radio coverage cells (also known as “sectors”) are arranged to provide static contiguous wireless coverage within a regional service area. Users of a mobile cellular network connect to the network via user equipment terminals (UEs), which attach wirelessly to one or more of the cells.\nThe network and the UEs support handover services, which provide continuous and uninterrupted user-perceived communication as users move between coverage cells. For example, the network may handover a UE from a first cell to a second cell when the UE moves from the first cell into the second cell.\nBecause of a random distribution of users and a varying demand for voice and data services, traffic loads vary from cell to cell. Adjoining or overlapping cells may have dramatically different loads over time.\nWhen a cell experiences more demand for resources than it can satisfy, the cell is “congested.” When the cell is congested, overloaded shared resources within the cell may be allocated on a fair basis among the UEs attached to the cell. However, each UE is allocated fewer resources when the cell is congested than when the cell is not congested. Those reduced allocations during congestion may be insufficient for certain high data throughput services, such as video services or timely large file delivery services, e.g., downloading email attachments. Congestion is therefore a condition to be avoided whenever possible.\nAccordingly, mobile networks, when possible, “load balance” cells by handing over UEs between serving cells. Load balancing prevents scenarios, for example, where one cell is operating in congestion, and another neighboring or overlapping cell is idle. Load balancing tends to maximize the usage of the aggregate capacity of the cellular network, improve the economic use of deployed infrastructure, and improve quality of experience for users of the mobile network.\nFor example, operators may trigger load balancing between cells whenever a numerical threshold difference exists in the aggregate traffic loads between adjoining or overlapping coverage cells that have one or more operating UEs in a common coverage area. For example, by forcing or biasing handover of UEs from a higher traffic load cell to a lower traffic load cell, the aggregate traffic load levels between the two cells are driven towards equilibrium, and general congestion may be lessened.\nExisting methods, however, currently lack novel methods of triggering and executing load balancing to more efficiently and equably reduce overall network congestion that may be aligned with Self-Organizing-Network (SON) principles."}
{"text": "The present invention pertains to heating, ventilating and air conditioning (HVAC) systems in general, and to an air handling unit arrangement in which a direct expansion coil is utilized.\nIn some buildings, typically high rises, it is common to use one or more small air handling units per floor. These systems have the advantages of being inexpensive to purchase and install and a self-contained system may be provided for each tenant. For example, each floor of a high-rise building may therefore have one or more small air handling units.\nSuch systems are characterized by recurring problems related to equipment failure and occupant discomfort. The recurring equipment problems can be identified as being related to icing of the expansion coil and cooling compressor seizure.\nThe occupant discomfort problems typically are associated with wide variations in temperature due to compressor cycling and excessive removal of moisture from the air."}
{"text": "Network carriers, also referred to sometimes as telecommunications operators or communications service providers, that run existing networks desire to optimize the network utilization for passing traffic, such as Internet Protocol (IP) traffic, over a physical network, e.g., across the Open Systems Interconnection (OSI) network layers one to five. The optimized traffic may include traffic for triple play services (e.g., video, voice, and/or data) and any type of bulk data. In existing networks, end-to-end services are typically set-up by Operational Support Systems (OSS) systems or provider specific network management service applications. Network carriers have suggested two different scenarios for optimizing network utilization and traffic: optimizing existing network services and enabling new/emerging network application services."}
{"text": "In an electronic device, two wires are usually connected to a plug connector and a receptacle connector, respectively, to facilitate an electrical connection of the two wires. The electrical connection of the two wires is achieved by mutual engagement between the plug connector and the receptacle connector.\nExisting plug connectors and receptacle connectors each comprise a housing and a connection terminal mounted within the housing. For ease of operation, the wires are generally electrically connected to the connection terminal in such a manner that a portion of the connection terminal pokes in a conductor portion of the wire. However, the existing plug connector and receptacle connector are used separately in pair. In case of connecting a plurality of wires, a plurality of plug connectors and a plurality of receptacle connectors are required. However, the plug connector and the receptacle connector in pair are arranged independently, which is inconvenient to be operated and managed."}
{"text": "This invention relates to medical devices that are resistant to bacterial growth or encrustation such as a) urinary catheters and more particularly to urinary catheters constructed of a material which enables the urinary catheters to inhibit urease and to prevent calcium and magnesium phosphate deposits on the catheters and b) contact lenses, more particularly to extended wear lenses coated with a suitable material to enable the lenses to resist the adherence of microorganisms.\nUrinary Catheters: Catheters are used in urological surgery and when other methods of managing urinary incontinence fail. These catheters are generally made from a nontoxic, non-irritating material which is flexible but which will withstand collapse. Ideally they should resist colonization by bacteria and encrustation by mineral deposits.\nThe most common materials used for the catheters are latex, plastic or silicone. All biomaterials, including the most recently tested polyurethane, polyurethane-carbon and silicone-carbon, become encrusted with mineral deposits to varying extents. The degree of encrustation formation on biomaterials exposed to urine is dependent upon the biomaterial, the length of urinary exposure, the presence of infection and the solute content of urine. Several independent studies (e.g. Hukins et al., 1989) of the scraped encrusted material have identified - two major salts, namely apatite Ca.sub.5 (PO.sub.4).sub.3 (OH) and struvite Mg NH.sub.4 PO.sub.4 6H.sub.2 O. A very small amount of brushite CaHPO.sub.4.2H.sub.2 O has also been identified. Apatite is precipitated from urine under alkaline conditions (pH=9.2) along with struvite. Brushite with a K.sub.sp of the order of 10.sup.-5 is not stable above pH=7; but is precipitated below 6.5.\nThe pH increase of urine occurs due to urease producing bacteria, that infect the urine. Urease catalyzes the hydrolysis of urea at an enormous rate (10.sup.14 times as fast as the spontaneous hydrolysis of urea which is not observable in neutral solution) producing ammonia, which in turn precipitates calcium phosphates.\nMicrobial heterogeneity and structural complexity of the biofilms observed on infected catheter materials by means of scanning and transmission electron microscopy reveal a variety of microflora. It appears that the adherent microcolonies of bacteria are much less susceptible than their planktonic counterparts to antibiotics because the extensive anionic matrix surrounding the cell appears to comprise an ion exchange barrier between charged antibiotic molecules and their cellular targets. Antibiotics should be reserved for symptomatic, febrile infections; otherwise their use produces a change in bacterial flora with the potential for producing resistance.\nBacteriurea in a catheterized patient indicates the urinary tract has become colonized or infected. Administration of anti-bacterials such as nitrofurantoin, methenamine, and naldixic acid are used in practice. Most of these are not very effective in alkaline urine, and to achieve a high concentration in urine, the patient has to be given enormous doses.\nCatheter associated bacteriurea results from ascending bacterial colonization within glycocalyx enclosed biofilm on the inside and/or outside surfaces of the catheter and drainage systems. Urinary tract catheters made from a biomaterial that inhibits bacterial adherence and thus retards upstream colonization of bacteria may reduce acquired urinary tract infection (UTI). While hydrophilic surfaces show reduced protein adhesion, antibacterial surfaces offer better protection. Either a controlled release or an antibacterial immobilized on the surface should offer continuous protection against bacterial colonization. A silver oxide catheter is described in a publication by Shaeffer, A. J., Story, K. O., Johnson, S. M. in \"Effect of silver oxide/Trichloroisocyanuric Acid Antimicrobial Urinary Drainage System on Catheter Associated Bactiurea\", J. Urol., 39, 60 (1988). A recent clinical study on silver catheters reported by Johnson, J. R., Roberts, P. L., Olsen, R. J., Moyer, K. A., and Stanni, W. E. in \"Silver Oxide Coated Catheters\", J. Infect. Diseases, 162, 1145 (1990) ) indicated the prevention of UTI among women not receiving antimicrobials.\nThe other problem associated with catheters left in place for long periods of time is the encrustation due to the formation of calcium hydroxyapatite and struvite. Long-term urethral catheterization is frequently necessary for patients with intractable urinary incontinence or retention. This procedure is employed in as many as 16-28% of patients in various chronic care facilities. More than 50% could suffer from blockage of their catheters. This can cause incontinence due to urinary bypassing of the catheter or acute pain and discomfort associated with urinary infection. This is particularly distressing for patients being cared for in the community where professional help is not immediately available. In addition, the coarse irregular surface may result in pain and physical trauma when the catheter is removed increasing the risk of infection.\nRemoval of deposits by acidic solution (citric acid-magnesium oxide) has been shown to dissolve encrustation in vitro and is now used extensively in practice. Frequent irrigation could result in the damage of the mucus in the bladder.\nThe patent literature abounds with patents on indwelling urinary catheter systems aimed at preventing urinary tract infection. These are mechanical devices with appropriate valve fittings to keep the drainage open at the same time preventing ascending infection. U.S. Pat. Nos. 4,946,449 issued to Richard Davis and 4,878,901 issued to Hans Ernst Sachse are two typical examples. However, these do not have any prevention methods for encrustation. In general, mechanical systems are complicated and do not provide adequate protection from infection. U.S. Pat. No. 4,932,948 issued to Kernes et al., discloses the use of a funnel shaped insert at the end of the urinary catheter that serves as a reservoir to antimicrobial agents. The antimicrobial agent is simply mixed with the ethylene/vinyl acetate polymer during the fabrication of the funnel. The catheter surface itself does not carry any antimicrobial nor is it capable of preventing encrustation. U.S. Pat. No. 4,579,554 issued to Jacob Glassman discloses a design that provides for irrigation of the catheter tube. Frequent irrigation, however, could result in the damage of the mucus in the bladder.\nU.S. Pat. No. 4,642,104 issued to Sakamoto et al., teaches the use of polymers carrying multicarboxyl, amino or sulfonic acid group capable of binding antibiotics through ion exchange. The ion exchange groups are chemically introduced into the molecules of the inside and outside wall of the urethral catheter by hydrolysis of certain functionalities that are present in the polymer that is coated on the surface of the catheter. Cationic antibiotics such as polymyxins or soap preparations such as benzalkonium chloride or benzethionium chloride, cyclohexidine or povidine-iodine remain on the surface due to the electrostatic binding.\nU.S. Pat. No. 4,950,256 issued to Luther et al., teaches the use of an intravascular catheter comprising a cannula for insertion into a vascular system of a patient. This catheter is coated with a hydrophilic polyurethane-polyene composition for binding antithrombogenic materials and cationic polymyxin antibiotics. Luther et al., show that the absorption of the polymyxins into the hydrophilic polyurethane can be controlled by varying the initial concentrations. The authors describe this as a time release intravascular catheter. The present invention differs from that of Luther et al., by virtue of the immobilization of polymyxin through a binding action with the dianionic diaminopropanol tetraacetic acid (DPTA) groups.\nContact Lenses: Infective ulcerative keratitis is one of the most severe hazards of hydrogel extended wear soft contact lenses (SCLs). The extent of bacterial adherence varies with the nature of the base materials, protein deposits and surface charges. While the mechanism of adhesion and colonization of bacteria on SCLs is not well understood, the influence of adsorbed proteins has been implicated. Recently it was reported (\"Pseudomonas attachment to low-water and high-water and non-ionic, new and rabbit worn soft contact lenses\", Brussel et al, Investigative Ophthalmology and Visual Science, 32, 657, 1991) that both new and worn SCLs can bind amounts of p. aeruginosa that could potentially produce bacterial keratitis.\nSCLs in the eye are readily coated with tear proteins, mucus and lipids (Tripathy et al., \"Lens Spoilage in Contact Lenses\", in The CLAO guide to basic and applied clinical practice\", Dabecies O. H. Jr. (Ed), Grune and Stratton Inc., Orlando Fla., 1984). P. aeruginosa is the responsible microorganism in up to two thirds of the reported cases of extended wear SCL associated infectious keratitis. The extended wear contact lens induces corneal hypoxia, microtrauma and tear film destruction that weaken corneal defense.\nIt is therefore a principal object of the present invention to provide a catheter for use in urological surgery or for managing urinary incontinence which inhibits the build up of encrustation on the surface of the catheter.\nAnother object of the present invention is to provide a catheter for use in urological surgery or for managing urinary incontinence which causes a lower risk of infection than known catheters.\nIt is yet another object of the present invention to provide a coating for contact lenses, particularly of the extended wear type, which inhibits the adhesion of harmful bacteria on the lens surface."}
{"text": "Locked Nucleic Acid (LNA) monomers and oligonucleotides were invented in 1997 (WO9914226) and have showed promising results as antisense drug candidates. However, there is no production method described in the literature for large scale LNA-phosphoramidite synthesis providing optimal yields of the LNA amidites necessary for efficient preparation of LNA based antisense drugs.\nOligonucleotide synthesis is typically performed using the phosphoramidite method, invented by Caruthers, U.S. Pat. No. 4,415,732, in 1980 and improved a couple of years later by Köster U.S. Pat. No. 4,725,677. LNA oligomers are being synthesised according to the phosphoramidite method, but so far the large scale supply of LNA monomers has been a problem due to slow reactions, side product formation during the reaction and reagents unstable at room temperature have been employed.\nThe mechanism of phosphoramidite activation and coupling in oligonucleotide syntheses has been studied in detail. Traditionally tetrazole has been used (Dahl et al. 1987, Nucleic Acid Research, 15, 1729-1743; Beaucage & Iver, 1992, Tetrahedron, 48, 2223-2311). The proposed mechanism of tetrazole is two step, first tetrazole protonates trivalent phosphorous followed by displacement of the N,N-diisopropylamine by the tetrazolide. This latter intermediate is very reactive with hydroxynucleophiles such as 5′-OH on nucleic acids. Therefore, tetrazole acts both as acid and as nucleophilic agent.\nVargeese et al. (Vargeese, C.; Carter, J.; Krivjansky, S.; Settle, A.; Kropp, E.; Peterson, K.; Pieken, W. Nucleic Acid Research, 1998, 26, 1046-1050; WO9816540;) have described an activator for the coupling of phosphoramidites to the 5′-hydroxyl group during oligonucleotide synthesis. The activator is 4,5-dicyanoimidazole (DCI) and its effectiveness is thought to be based on its nucleophilicity. It was shown that DCI significantly increased the yield in phosphoramidite oligomerisation.\nVargeese et al. has described a method of small scale DNA phosphoramidite thymine monomer synthesis using DCI, but the method they used provided a moderate yield (75%) after several re-precipitations of the amidite. Kittaka et al. (Kittaka A., Kuze T.; Amano M.; Tanaka H.; Miyasaka T.; Hirose K.; Yoshida T.; Sarai A.; Yasukawa T. and Ishii S. Nucleosides & Nucleotides 18, 2769-2783, 1999) have also used DCI for amidite synthesis but they reported an even lower yield of the product (65%). Kittaka et al. (Kittaka A.; Horii C.; Kuze T.; Asakura T.; Ito K.; Nakamura K. T.; Miyasaka T. and Inoue J., Synthetic letters, S1, 869-872, 1999) have also made a derivatised uridine phosphoramidite by 3′-O phosphitylation using DCI (0.7 eq), but without describing the reaction parameters further. The yield was reported to be 92%. Generally, phosphitylation of uracil and thymine provide the highest yields because these nucleobases usually do not interfere with the phosphitylation reagents.\nAccordingly, there is nothing in the literature describing that DCI has been used as an activator for LNA phosphoramidite synthesis.\nThe previously reported oxy-β-D-ribo-LNA (FIG. 2) method of synthesis used 2-cyanoethyl-N,N-diisopropylchlorophosphoramidite and diisopropylethylamine (Koshkin, A. A.; Singh, S. K.; Nielsen, P.; Rajwanshi, V. K.; Kumar, R.; Meldgaard, M.; Olsen, C. E.; Wengel, J. Tet. 1998, 54, 3607-3630). The reported yields were low (LNA-T 70%, LNA-U 58%, LNA-G 64%, LNA-A 73% and no yield for LNA-C was reported). The previously reported method for synthesising the oxy-α-L-ribo thymine phosphoramidite (FIG. 3) also utilized 2-cyanoethyl-N,N-diisopropylchlorophosphoramidite and diisopropylethylamine and gave a yield of 60% (Hakansson, A. E.; Koshkin, A. A.; Sorensen M. D.; Wengel J. J.Org.Chem. 2000, 65, 5161-5166)."}
{"text": "In order to determine the stability, reactivity and compatibility of two or more materials, such as explosives, paints or adhesives, it is necessary to monitor the change of pressure of these materials over an extended period of time.\nCurrent standard techniques employed to monitor such change of pressure utilize a mercury manometer connected to a sample tube. The materials are placed in a test tube which is fitted with a ground glass ball joint. A manometer tube is attached thereto via a mating ground glass ball joint. The reservoir on the other end of the manometer is filled with mercury. A vacuum source is attached to the tip of the manometer to permit evacuation of the interior of the assembly. After evacuation, the assembly is placed in a holding fixture with the test tube immersed in a heated bath at a preferred temperature of 100 degrees Centigrade.\nThe distance from the top of the mercury in the reservoir to the top of the column in the manometer capillary is measured and recorded, along with the ambient temperature and ambient pressure. At the end of a specified time, usually 40 hours, the mercury column height and ambient temperature and pressure are measured and recorded again. The change in internal gas volume is then calculated, corrected for standard temperature and pressure.\nThe data obtained by this technique provides only the change in internal gas volume over the duration of the test. The fragility of the manometer makes it a practical impossibility to measure the column height with any frequency. Moreover, the technique requires the handling of mercury in large quantities. Precision calibrated glassware needs to be repeatedly handled, including cleansing of melted energetic materials adhering to the test tube. Operating costs are high due to repeated cleaning requirements, the high cost of replacing breakable precision glassware, and the need to perform extensive and complicated calibration procedures on new glassware components."}
{"text": "This invention relates generally to non-volatile semiconductor memories, their operation and, in particular, to the operation of memory systems that include multiple dies.\nThere are many commercially successful non-volatile memory products available today, which use an array of flash EEPROM cells. An example of a flash memory system is shown in FIG. 1, in which a memory cell array 1 is formed on a memory chip 12, along with various peripheral circuits such as column control circuits 2, row control circuits 3, data input/output circuits 6, etc.\nOne popular flash EEPROM architecture utilizes a NAND array, wherein a large number of strings of memory cells are connected through one or more select transistors between individual bit lines and a reference potential. FIGS. 2A-2B illustrate an example of a planar NAND flash memory array. In other examples, NAND strings extend in a vertical direction in what may be referred to as three dimensional (3D) memory.\nFlash memory is generally arranged in blocks, with a block being the unit of erase. FIG. 3A illustrates blocks in a memory die that are arranged in two planes. Blocks in a plane share certain circuits so that only one block in a plane is accessed at a time. Multiple planes allow multiple blocks in a die to be accessed in parallel.\nMultiple dies may be connected to a memory controller by a memory bus as shown in FIG. 3B. The memory controller receives data and distributes it to the dies. While such an arrangement may allow a high degree of parallelism in some conditions, delays may occur when one or more dies become busy and thus may affect write speed."}
{"text": "1. Field\nEmbodiments of the invention relate to garbage collection in string tables having strings that are no longer referenced.\n2. Description of the Related Art\nString interning is a commonly used technique in which a string table is used to reduce the space required for strings and to make string comparisons faster. String interning may be described as a technique of storing one copy of each distinct string value. Each string is replaced with an integer (or pointer) which is associated with that unique string. String comparisons are reduced to single integer comparisons. The space used by many copies of the same string is reduced to the space used by just the single copy in the string table. An interned string table may be described as a string table for which string interning is performed.\nThus, with some techniques, strings may be interned into a string table such that interned references to the same string (or string value) reference the same in-memory instance. That is, the reference refers to the string. In particular, if the integer or pointer for the interned string is stored, that integer or pointer is called a reference. Also, a reference may be described as a handle stored somewhere, but the reference refers to the string. That is, computer code may make use of the string, so the string can not be deleted. While string tables may provide memory and performance advantages, the string tables may overfill with strings (i.e., the string table may run out of space) or the integer type (which could be a 16-bit short word) may run out of codes if there are too many unique strings. This may lead to a denial of service vulnerability if the strings which are interned arrive over a network because an attacker need only send many unique strings to the server until the string table is full or out of codes. As another example, for an XML document having a large number of names, storing the names in a string table may cause the string table to overflow. If used carelessly, a string-table overflow may lead to poor table lookup performance, may exhaust a fixed-size handle pool, or may even exhaust system memory, resulting in performance and/or stability problems.\nThese problems can be solved with some form of garbage collection of the string table, in which strings which are no longer referenced are discarded.\nIn environments with manual memory management or when stricter control over the string table size or lifecycle is desired, recovery of unused strings is complex and specialized string table garbage collection may be required. In particular, in the absence of a system-wide module to determine reachability of the string table entries, the string table has to maintain its own data on the strings in use. This can be done with reference counts, but that will incur penalties in performance, and possibly memory, since the string table users will constantly be updating reference counts on the strings, and since such reference-counting mechanisms often depend on heavier-weight structures to deal with the reference-count update. In a system where interned strings or string handles are used extensively and passed around among various modules, the overhead of reference count updates may reduce the performance advantage that the string table provides.\nIn environments (i.e., platforms, languages, and/or frameworks) that provide native garbage collection, especially those with support for weak references, a well-designed string table may allow its entries to be freed by the system garbage collector, allowing free use of the string table without risking overflow. However, even in environments where native garbage collection is available, a performance-sensitive or memory-sensitive application may want to use lighter-weight references to the strings, such as small integers, or may wish to have greater control over the size of string-table or the life-cycle of the entries. For example, using a smaller fixed-size table may prove to be easier to tune for performance, but may become exhausted before the system garbage collector decides to free the entries. Similarly, using small-integer handles for strings may allow optimization of other handle-keyed tables throughout the system.\nWith concurrent garbage collection, other processing does not have to stop while garbage collection is underway. However, with existing garbage collection techniques, none are completely concurrent. For example, some existing garbage collection techniques lock the root set (i.e., the stack) during part of the garbage collection.\nIn a distributed environment having multiple nodes, memory addresses can not be used as the integer keys because the various nodes do not share the same main memory or virtual memory address spaces. Thus, the nodes can not allocate an entry in the string table for a new string without first synchronizing and checking that no other nodes have already allocated a string at that memory location. The overhead of synchronizing these memory addresses alone would be high-cost and high-complexity. Some systems run completely separate processes on the various nodes and pass the full strings between nodes. This has the drawback that the optimized data structures used on one node can't be passed between nodes.\nThus, there is a need for improved garbage collection techniques."}
{"text": "1. Field of the Invention\nThe present invention relates to a resin composite structure obtained by integrally bonding fiber-reinforced thermoplastic resin composite members containing discontinuous fibers. The present invention also relates to a method for producing such a resin composite structure.\n2. Description of the Related Art\nA fiber-reinforced thermoplastic resin composite (hereinafter referred to simply as a resin composite) is a material containing fibers such as glass or carbon fibers impregnated with a thermoplastic resin. In recent years, such a resin composite has been studied as a material for a component of a car body (such as a floor or a monocoque component).\nIn general, the car body component has a large size. Therefore, in order to produce the entire component in one molding process, a large scale molding apparatus is needed to apply a large load during the molding process. Such a large scale molding apparatus is expensive and consumes a great amount of power during the molding process. Therefore, this type of molding apparatus is disadvantageous for mass production in terms of its high facility cost and high running cost.\nA large component can be obtained by connecting a plurality of members that have been prepared separately. However, in the case that such separate members are connected mechanically using nuts and bolts or rivets (i.e., so-called metal fasteners), a boring operation and a fastening operation are required, thereby resulting in high operating costs. Furthermore, the weight of the component is increased due to the metal fasteners. In addition, in this case, electrolytic corrosion occurs at a contact portion between the reinforcing fibers and the metal fasteners. Further, it may be necessary to take countermeasures in response to different coefficients of thermal expansion of the materials.\nThe members may be connected by other known processes such as thermal welding, vibration welding, ultrasonic welding, or laser welding. However, such processes do not produce a sufficient welding strength.\nIn a method proposed in Japanese Laid-Open Patent Publication No. 11-090986, the members first are heated and melted such that the reinforcing fiber is partially exposed on the surface of the connection, and thereafter, the melted members are press-molded and welded. Japanese Laid-Open Patent Publication No. 11-090986 describes that the connection obtained by this welding method has an excellent appearance and exhibits a high welding strength."}
{"text": "1. Technical Field\nThe present disclosure relates to an input inducing device and an inducing keyboard.\n2. Description of Related Art\nThe commonly used input device such as a keyboard is operated by pressing buttons. However, the buttons of the keyboard are easily damaged because of the repeated and long time pressing."}
{"text": "1. Field of the Invention\nThe invention relates generally to devices, methods, and systems for securing data bearing media. Specifically, the invention relates to devices, methods, and systems for preventing unauthorized access to data bearing media.\n2. Description of the Related Art\nData storage devices, such as disk drives, are currently used in a variety of machines including computers, car stereos, vending machines, media players, and automated teller machines (ATMs). In the foreseeable future, data storage devices will become increasingly portable and will be introduced into additional machines and environments due to their ever-increasing storage capacity, shrinking footprint, and decreasing price.\nData storage devices often contain large amounts of sensitive data such as financial information, personal information, business plans, and more. As data storage devices become increasingly pervasive, interchangeable, and portable, the ability to secure data storage devices, particularly portable devices, is becoming increasingly important.\nOne of many methods for breaching a data storage device is to remove the data storage device from a machine and access the data storage device in another location with a machine or computer system configured for that purpose. For example, an unauthorized party may attempt to access the data storage device of an ATM by physically removing the data storage device from the ATM, transporting it to a remote location, and connecting it to a computer system designed to generate a large number of passwords and associated access attempts in a short period of time. Consequently, the unauthorized party may gain access to the valuable information stored upon the stolen data storage device.\nOver the years, a variety of techniques have been developed that attempt to protect data storage devices from unwanted intrusion. For example, some data storage devices prompt the user for multiple passwords in order to decrease the probability of successful intrusion. Other storage devices impose a delay before responding to a password in order to reduce the rate at which passwords may be entered. While useful in reducing the likelihood of intrusion, such methods may be overcome by a persistent intruder.\nWhile invulnerable protection from intrusion is unattainable, what is needed is a data storage device that further reduces the probability of successful intrusion by an unauthorized party. More specifically, what is needed is a data storage device that prevents unauthorized access by imposing timing requirements on password related data and an associated computing device that provides password-related data with the prescribed timing requirements imposed thereon."}
{"text": "A production method of an optical display device mounted to a conventional liquid crystal display device is conceptually shown in FIG. 11. First, in an optical filmmaker, a step produces a long (web-like) sheet material having an optical film as a material roll (#1). The concrete production step is a known production step, and a description thereof will not be given. As the “long (web-like) sheet material”, for example, there are a polarizing plate material, a retardation plate material, a laminated film material of the polarizing plate and the retardation plate, and the like which are used in a liquid crystal display device. Next, the material roll is slit to a predetermined size (a size in accordance with a size of the optical display unit) (#2). Next, the slit long material is cut to a fixed size in conformity to a size of the optical display unit (#3). Next, a step inspects an outer appearance of a piece of sheet material (an optical film) cut to the fixed size (#4). As the inspecting method, for example, there can be listed up a defect inspection in accordance with a visual observation, and an inspection using a known defect inspection apparatus. The defect means, for example, a dirty in a front face or an internal portion, a scratch, a special twisted defect like a hitting mark generated by biting a contaminant (which may be called as a knick), an air bubble, a contaminant or the like. Next, a step inspects a finished product (#5). The finished product inspection is an inspection in accordance with a quality standard having a severer non-defective determination than the outer appearance inspection. Next, a step works end faces in four sides of the sheet material of the piece of sheet material (#6). This step is carried out for preventing an adhesive or the like from running over from the end faces during transport. Next, a step cleanly packages the piece of sheet material under a clean room environment (#7). Next, a step packages for transport (a transport package) (#8). The piece of sheet material is produced as mentioned above, and is transported to a panel processing manufacturer.\nIn the panel processing manufacturer, a step dismounts the package of the piece of sheet material transported (#11). Next, a step inspects an outer appearance for inspecting the scratch, the dirt and the like generated at a time of transporting or dismounting the package (#12). The piece of sheet material which is determined as the non-defective by the inspection is fed to the next step. There is a case that the outer appearance inspection is omitted. An optical display unit (for example, a glass substrate unit in which a liquid crystal cell is enclosed) to which the piece of sheet material is bonded is previously produced, and the optical display unit is cleaned before the bonding step (#13).\nA step bonds the piece of sheet material and the optical display unit (#14). A release film is peeled off from the piece of sheet material while leaving a pressure-sensitive adhesive layer, and it is bonded to one face of the optical display unit by using the pressure-sensitive adhesive layer as a bonding face. Further, it can be bonded to the other face of the optical display unit in the same manner. In the case of bonding to both the faces, the structure may be made such that the optical films having the same construction are bonded to the faces of the optical display unit, or the structure may be made such that the optical films having different constructions are bonded thereto. Next, a step carries out an inspection of the optical display device in the state in which the optical film is bonded and a defect inspection (#15). The optical display device which is determined as the non-defective in this inspection is fed to a mounting process (#16). On the other hand, a reworking process is applied to the optical display device which is determined as a defective (#17). In the reworking process, the optical film is peeled off from the optical display unit. The optical film is newly bonded to the reworked optical display unit (#14).\nIn the production step mentioned above, since the optical film maker and the panel processing manufacturer exist in the separate places, the end face working, the packaging of the piece of sheet material, the package dismounting and the like are particularly necessary steps. However, there are a problem of a production cost increase caused by multiple steps, a problem of the scratch, the dust, the dirt and the like generated by the multiple steps and the transport, a necessity of the inspection step caused thereby, and a problem that it is necessary to store and manage many kinds of sheet materials as a stock.\nAs a method of solving the problems, there has been proposed Japanese Patent Application Laid-Open (JP-A) No. 2007-140046. In accordance with this invention, the structure is provided with a supply portion pulling out and supplying a long sheet material from a material roll around which the long sheet material having an optical film corresponding to a member of an optical display device is wound, a detection portion detecting a defect of the long sheet material pulled out by the supply portion, a cutting work portion cutting the long sheet material based on a result of detection of the detection portion and working to an individual sheet material, a transfer portion transferring the sheet material cut by the cutting work portion for a bonding work, and a bonding work portion bonding the sheet material transferred by the transfer portion and an optical display unit corresponding to a member of an optical display device, and these portion are arranged on a continuous production line. In the structure mentioned above, it is possible to directly cut the long sheet material having the optical film into a desired size, and to bond the cut sheet material to the optical display unit. Accordingly, it is possible to directly package the long sheet material wound around the material roll so as to deliver, in place of the conventional step which stamps the long sheet material, tightly packages the stamped sheet material, and delivers to the panel processing manufacturer."}
{"text": "Existing 802 technology (802.11 WLANs, 802.15 wireless personal area networks (WPANs), etc.) traditionally does not have to support emergency calls like cellular does. For cellular, support of emergency calls often resulted from regulatory requirements imposed on the technology and is therefore widely implemented in most of today's deployed wireless cellular networks and handsets. Support for emergency calls involves many aspects across all communication layers, especially signaling support and mandated procedures, which are non-existent for 802.11 and 802.15 technologies. With the advent of Voice over Internet Protocol (VoIP) in WLANs and increased everyday usage of WLANs, support for emergency calls in WLANs will become necessary.\nEven “fixed” VoIP phone service offerings for the residential market have limited emergency call support. Number location information cannot always be tracked by a dispatcher in a public safety answering point (PSAP), call back is not always possible, and address registration may be required upon purchase of the equipment. When the VoIP phone is moved to a new location, the emergency call will still be sent based on the registered address location. The registered address can be changed in principle, but delays are at least on the order of days or weeks in updating the information at the PSAP. In addition, some users might not update their registration information in a timely manner, if at all.\nThis situation worsens with more mobility as enabled by VoIP phones using WLANs. WLAN based VoIP phones can work from any location and the user can be expected to roam seamlessly between locations, such as from an office to a home to public hotspots, etc.\nCertain 802.11-specific issues exist, including radio access, access point (AP) location, caller location, and emergency call admission. In regard to radio access, no priorities for emergency calls currently exist in the 802.11 standards, and there is no means to distinguish an emergency call from a regular call for the WLAN access network. The location of an AP or a STA is currently unknown to the network in a non-proprietary manner, even if for example, the AP's identification can easily be determined. It is also not currently possible to map the caller's location in a non-proprietary manner.\nIn regard to admission, a tightly managed WLAN may prevent emergency callers from establishing an emergency call if the caller is not authorized to enter the network. The normal connection procedure between a STA and an AP requires the STA to send an Association request, followed by negotiation with the AP prior to associating the STA to the AP. If the STA is unable to indicate that it is making an emergency call, it would have to go through the entire association procedure to determine if it could be admitted. As an example of this type of difficulty, if a STA does not have the proper password or authentication credentials to access the system (if the AP is configured to require passwords or require authentication credentials, as might exist for example with a private hotspot or enterprise/office WLAN), the AP will bluntly refuse the STA's association request. But even if the STA has the proper password or authentication credentials, the AP could still refuse admittance to the network based on its configured maximum capacity for voice users. In this case, the right decision for the AP would be to admit this new emergency call (at the highest priority) and to discontinue another existing voice call. Because the AP currently lacks means to make this distinction in the first place, such a feature cannot be implemented with existing state-of-the-art WLAN technology. Contrast this with the operation of a cellular system, in which any device can make an emergency call, even a device without a SIM card."}
{"text": "The present embodiments relate to aggregation operations. More specifically, the embodiments relate to translating a query involving a distinct aggregate(s) into a query that does not involve a distinct aggregate.\nData aggregation is a process in which information is gathered and expressed in a summary form. Submitted data may be processed by an aggregate which is used to compute user specified operations over user specified data ranges. Examples of such aggregates include, “AVG” which functions to return the average of values defined in the expression, “BASE” which returns the total number of cases includes in the expression, “COUNT” which functions to return a count of cases in a column, “DISTINCT” which functions to return unique values, “MIN” which functions to return the lowest value, “MAX” which functions to return the highest value, “STDEV” which functions to return the standard deviation of the values defined in the expression, and “SUM” which functions to return the sum of the values defined in the submitted expression. In addition to these example aggregates, a distinct aggregate is defined as an aggregation performed on the results of a DISTINCT operation, e.g. SUM (DISTINCT (value(s))).\nOne popular database language that may be used to perform aggregations is Structured Query Language (SQL). An SQL statement is used to query or manipulate data in data storage. For example, a SELECT, UPDATE, INSERT or DELETE statement can be used with the SQL statement. A query using a SELECT statement may contain the following clauses, SELECT, FROM, WHERE, GROUP BY, HAVING, ORDER BY, and JOIN. The SELECT clause is used to specify the data of interest. The FROM clause indicates a source for the data. The WHERE clause indicates a condition for filtering the data. The GROUP BY clause allows organization of data into subsets. The HAVING clause places a condition on membership in a subset. The ORDER BY clause organizes the resultant data. The JOIN clause matches a row of one table with a row of another on the basis of a join condition. The query presented in the SQL statement is not limited to a single clause.\nAny query under workspaces can suffer performance issues (e.g. become time consuming and memory intensive). In order to compute a query containing a distinct aggregate, a calculation of the aggregate's corresponding distinct values in the database must be performed. Queries containing distinct aggregates are particularly susceptible to suffering performance issues especially when the query contains multiple distinct aggregates since each aggregate has to be computed separately for values belonging to each category defined in the aggregate. Further, the performance issues with a query become exacerbated when performed in distributed environments, such as the cloud, due to the structure of the data storage and plurality of data storage options. Accordingly, it is desirable to provide a system and method for executing distinct aggregate operations in a manner that is less memory intensive and time consuming than current techniques."}
{"text": "This invention relates generally to concrete slab forming machines, and has particular reference to an improved machine for forming high quality concrete slabs by a continuous process without the use of conventional casting forms.\nForming slabs by a continuous process without forms requires a concrete mix of a consistency such that the concrete section that is formed will be self-supporting before the concrete has hardened. In concrete technology, such concrete mixes are defined as \"zero slump\" and \"low slump\" (up to 3 cm .+-. slump) mixes.\n\"Zero slump\" (dry state) mixes are highly abrasive and cause excessive wear on the forming machinery. An even more serious disadvantage, however, is the lack of cement-sand paste in such mixes which prevents the achievement of a good bond between the concrete and any reinforcing steel there may be.\n\"Low slump\" (plastic state) mixes produce structurally sound concrete with far less wear on the machinery and good reinforcing steel-concrete bonding. However, with such mixes there may be some sagging or slumping in the freshly formed section resulting in an inferior product.\nBetween these extremes of workable consistency limits, there is a range where the concrete mixes have optimum characteristics and mixes in this range are defined as \"zero slump saturated\" mixes. These mixes have a practical water variance flexibility and combine the advantages of the \"zero slump\" stiffness after the forming stage (inside and outside the machinery) and \"low slump\" workability during the forming stage.\nAt the present time, concrete slab forming machines can be broadly classified as extrusion type and slip form type machines. The extrusion type machines exert a force on the concrete mix to extrude it through the forming part of the machine. When the extruded concrete section reaches a certain resistance (weight), the reaction of the extrusion force propels the forming machine ahead. Within the extrusion machine classification there are screw type and reciprocating type machines.\nScrew type extrusion machines can use only \"zero slump\" concrete mixes because the mix must be stiff to transmit the longitudinal extrusion force from the cylindrical envelopes of the augers to the rectangular shape of the concrete section. If water were added to the mix to improve its bond with reinforcing steel, the mix would very possibly plug the augers. To achieve the necessary cement-sand paste for good bonding when using this type of machine, water from the floor is vibrated in the bottom of the concrete section. This, however, results in different densities between the top and bottom of the slab which causes excessive camber/deflection problems.\nReciprocating type extrusion machines can use concrete mixes from \"zero slump\" through \"low slump\" consistencies. However, because of the stop and go that is necessitated by the feed and extrude cycles, these machines have the disadvantage of relatively low production speed. U.S. Pat. No. 3,159,897 to Ellis discloses a screw type machine, while U.S. Pat. Nos. 3,143,718 and 3,143,782, both to Kalns, disclose reciprocating type machines.\nSlip form type slab forming machines, unlike the extrusion type, do not apply any force to extrude, or to pack the concrete mix into the forming part of the machine. These machines are propelled by external means such as a power wheel or winch and the concrete mix is moved by gravity and vibration into the forming part of the machine. To enable the mix to flow, and to minimize the friction between the forms and concrete, the mix must be in the wet range (2 cm .+-. slump) of the low slump consistency resulting in an inferior product. In addition to this disadvantage, the gravity feed frequently is not uniform due to the practical limitations in controlling the concrete mix. Water and aggregate changes influence the flow characteristics of the wet concrete mix and thus impair the uniformity of the product. Concrete mixes of even medium consistency (1 cm .+-. slump) cannot be used in the slip form machinery because of the lack of a force to compensate for the friction between the forms and concrete mix, i.e. such mixes plug up the machinery.\nSlip form type slab forming machines are disclosed in U.S. Pat. Nos. 3,401,438 to Dennis; 3,583,046 to Dickinson; 3,608,011 to Jones and 3,647,308 to Yost."}
{"text": "1. Field of the Invention\nThe present invention relates to a novel thermal printhead and, more particularly, to a thermal printhead having heat generating resistors made of a Cr-Si-SiO alloy.\n2. Description of the Prior Art\nHitherto, the heat generating resistor for use in thermal printhead has been produced either from (1) a Cr-Si alloy as disclosed in the specification of U.S. Pat. No. 4,343,986 or (2) a Cr-SiO alloy as disclosed in IBM Technical Disclosure Bulletin Vol. 23, No. 2, July 1980.\nThe heat generating resistor made of Cr-Si alloy, however, suffers from the following disadvantages (a) and (b), while the heat generating resistor made of Cr-SiO alloy suffers from the following disadvantages (c) and (d) in addition to the disadvantages (a) and (b) common to the heat generating resistor made of Cr-Si alloy.\n(Disadvantages of Heat Generating Resistor Made of Cr-Si Alloy)\n(a) The maximum electric power density is as small as several W/mm.sup.2.\n(b) Resistance value is undesirably changed when the heat generating resistor is used for a long time at a high temperature. For instance, 3 to 7% change in resistance value is caused as compared with the initial resistance value (resistance value after a heat treatment up to high temperature, for instance 400.degree. C., following the film formation by sputtering) after 600 hours operation at 200.degree. C.\n(Disadvantage of Heat Generating Resistor Made of Cr-SiO Alloy)\nThe following disadvantages (c) and (d) are encountered besides the above-mentioned disadvantages (a) and (b).\n(c) Only negative temperature coefficient is available so that the use is impractical for heat treatment resistor application.\n(d) Etching speed is too low.\nUnder this circumstance, there is an increasing demand for resistor material having improved performance as the material of, for example, thin film thermal printhead."}
{"text": "1. Field of the Invention\nThis invention is directed to an improved method for preparing alkenyl-substituted succinic anhydrides.\nAlkenyl-substituted succinic anhydrides are valuable intermediates for the production of products which have a wide range of uses. For example, the acids and the anhydrides are useful as rust inhibiting agents in mineral and synthetic lubricating oils and in paint and other protective coating compositions. The reaction products of alkenylsuccinic anhydrides with amines results in the production of alkenyl-substituted succinamic acids and succinimides and provides products which have valuable friction modifying and detergent-dispersant properties for lubricating oil compositions.\nThe alkenylation of maleic anhydride to produce an alkenyl-substituted succinic anhydride particularly with the higher molecular weight olefin polymers is a relatively slow, inefficient process. The art describes many processes for improving the reaction between an olefin polymer and maleic anhydride. One such process involves halogenating the olefin polymer prior to reacting with it maleic anhydride. Many processes for improving this reaction are based on the use of a catalyst promoter. Thus, processes have been proposed which involve the use of a halogenated organic or inorganic catalyst such as a chlorinated or brominated organic or inorganic compound. Another type of promoted process involves the use of a free radical initiator in the process.\n2. Description of the Prior Art\nU.S. Pat. No. 3,231,587 discloses a process for reacting an aliphatic polymer of a lower metal olefin and maleic anhydride by contacting the reactants at a temperature of above about 140.degree. C. in the presence of about 1 mole of chlorine for each mole of maleic anhydride.\nU.S. Pat. No. 3,927,041 discloses a method for preparing an alkenylsuccinic anhydride by reacting a viscous polybutene with an unsaturated aliphatic dicarboxylic acid anhydride in the presence of 1,3-dibromo-5,5-dialkyl-substituted hydantoin."}
{"text": "This section is intended to introduce the reader to various aspects of art, which may be related to various aspects of the present invention that are described or claimed below. This discussion is believed to be helpful in providing the reader with background information to facilitate a better understanding of the various aspects of the present invention. Accordingly, it should be understood that these statements are to be read in this light, and not as admissions of prior art.\nComputer systems often include a circuit board, a chassis, and electronic components. The circuit board is generally a motherboard, a system board, or a main board, and the electronic components often include a power supply, a sound card, a video card, a backplane, or the like. Typically, the electronic components connect to the circuit board, and the chassis supports the circuit board.\nUnfortunately, if the circuit board is not adequately secured to the chassis, the circuit board and electronic components may move and/or disconnect from one another. Additionally, if the circuit board is not adequately secured to the chassis, it may be difficult to reliably connect the circuit board to the electronic components. An unsecured circuit board may shift out of alignment with the electronic components and potentially render a connection between the two unreliable."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for driving a bistable display.\nThe present invention further relates to an apparatus for driving a bistable display.\nThe present invention further relates to a system comprising a bistable display and an apparatus for driving the same.\n2. Related Art\nMultistable displays, such as electrophoretic displays, have a plurality of pixels, which may be settable with a first operating luminance level, a second operating luminance level and an intermediate operating luminance level. Electrowetting based displays are another example of a multistable display technology. Also LCD based displays have been developed having a multistable behavior. Typically, multistable displays are reflection type displays. Accordingly the luminance level is determined by a reflection level. Alternatively, a transmission type multistable display may be displayed, wherein the luminance level is determined by a transmission level. Conventionally, multistable displays are denoted as “bistable displays”. This denotation will be used throughout the description. In the following the wording “luminance level” will also be briefly denoted as “luminance”.\nUsually, the first operating luminance level relates to “white”, the second operating luminance level relates to “black” and the intermediate operating luminance level relates to “grey”. In order to change image content on an electrophoretic display, new image information is written for a certain amount of time, for example during a period of 300 ms-600 ms. The refresh rate of the active-matrix is usually higher (for example 20 ms frame time for a 50 Hz display and 10 ms frame time for a 100 Hz display). Changing pixels of such display from black to white, for example, requires the pixel capacitors to be charged to a suitable control voltage for 200 ms to 300 ms, in the case where a pulse-width modulation principle is used. During this time the white particles drift towards the top (common) electrode, while the black particles drift towards the bottom electrode, for example an active-matrix back plane. Nevertheless, in order to rule out effects of earlier states of the display updating to an accurately defined new state requires an update time that is about three times longer, e.g. in the range of 600 to 900 ms. Switching to black requires a control voltage of a different polarity, and applying substantially 0 V on the pixel substantially preserves its condition. Addressing such electrophoretic display for a short time with a certain voltage will result in a situation wherein a mixture of white and black particles is visible. Alternatively, electrophoretic displays exist that use only one type of particle. Therein the perceived grey value is determined by the position of the particles with respect to the electrodes. Because the particles are very small human eyes integrate various ratios of black and white particles to shades/levels of grey. Such condition is regarded as an intermediate reflection level.\nBistable displays may have an infinite number of microstates depending on the momentaneous position and velocity of the particles that determine the luminance of the pixel. However, for practical purposes it will be presumed that the state of the pixel is one of a predetermined number of states that corresponds to a respective one of that predetermined number of grey values that is controlled by the apparatus for driving the display.\nWO 2009/078711 describes a method and apparatus for controlling an electronic display having a plurality of pixels settable in a plurality of reflection levels comprising a first level, a second level and a plurality of intermediate levels. The intermediate levels form a substantially equidistant partition of a dynamic range between the first level and the second level. The method comprises the step of setting the pixels to a preparatory intermediate level immediately prior to setting the pixels in a desired level selectable from said plurality of levels. The preparatory intermediate level can be selected from two or more levels. Subsequently, pulse width modulation is used to set the pixels in said desired level starting from the selected preparatory level.\nPixels of the known electrophoretic display have a limited bit depth. For example, a 4 bit pixel has 24=16 grey levels. In order to enable 32 levels (distinct shades) the pixels have to be controlled with a 5-bit driving scheme. For the known electrophoretic display for an equidistant partition of a full dynamic range of a pixel (e.g., between lightest to darkest shades), increasing the bit depth could require increasing the frame rate. Increasing the frame rate generally increases power consumption and potentially leads to a shorter product lifetime. Also, increasing the bit depth requires a higher accuracy and robustness of the method to control the display used to obtain the equidistant partitioning of the dynamic range."}
{"text": "The present disclosure relates generally to the field of heat pump systems using a vapor compression cycle, and more specifically to use of such systems in heating, ventilation, air conditioning, refrigeration, fluid heating and chilling.\nVarious systems are available for heating, ventilation, air conditioning, refrigeration, fluid heating and chilling. Such systems can be dedicated to heating or to cooling, or, for example, in the case of heat pump systems, the direction of a refrigerant flow can be reversed through heat exchangers in forced-air systems to either allow for absorption of heat from a space for the cooling such space or for absorption heat from the outdoors for the heating of such space. In this type of arrangement, forced-air flows over the heat exchanger through ductwork to such space.\nHeat pump systems may also be used in ductless systems, including direct expansion systems, where a refrigerant heat exchanger is typically in the space to be heated and/or cooled. Variable refrigerant flow systems are examples of direct expansion systems and may offer benefits in certain applications, including improved energy efficiency. Such systems may use one or more condensing units and provide refrigerant to one or more evaporator units in a ductless manner.\nIn certain situations, conventional direct expansion systems may involve safety issues in that because one or more of the heat exchangers is in the space being heated and/or cooled, in the event of a refrigerant leak, the refrigerant may leak into such space. Certain refrigerant gas is heavier than air and can displace oxygen in a room or space, and in extreme situations, could displace a sufficient amount of oxygen from a space such that a person could succumb to suffocation. As certain refrigerants cannot generally be detected by a person through sight, smell, or otherwise, the severity of such a refrigerant leak could become grave.\nTherefore, there is a need for heating and chilling systems for a fluid, such as water, which provide increased efficiency and which also lessen the risk of injury to persons in the event of a refrigerant leak."}
{"text": "The present novel concept broadly relates to the art of distance measurement and, more particularly, to a system and method for indicating the distance between associated structural members using electromagnetic wave modulation.\nThe subject system and method are amenable to broad use in a wide variety of applications and environments. One example of a suitable application is the use of the subject system and method on and with an associated fluid suspension member, such as an air spring of a vehicle, for example. The subject system and method will be discussed in detail hereinafter with specific reference to use on such an associated fluid suspension member. However, it is to be specifically understood that the subject system and method are capable of broader application and are not intended to be limited to the specific examples shown and discussed herein, which are merely examples of suitable applications.\nA variety of well known and commonly used devices and arrangements have been and are currently used to monitor the relative position of one structural member to another. For example, mechanical linkage sensors that include one or more linkage members are often used to connect between adjacent structural members, such as a suspension component of a vehicle and the corresponding frame or body of the same. The linkage members typically act through a variable resistor or other suitable component that changes in response to the movement of the linkage. An electronic control unit (ECU) or other suitable device then determines the relative position of one structural member to the other based upon a corresponding change in voltage across the variable resistor or a corresponding change in current through the resistor.\nUnfortunately, such arrangements have a number of problems and/or disadvantages that are commonly associated with their continued use. One problem with the use of mechanical linkages, particularly those used in association with the suspension system of a vehicle, is that the linkages are frequently subjected to physical impacts, such as may be caused by debris from a roadway, for example. This can result in the linkage being significantly damaged or broken, such that the device no longer operates properly, if it operates at all.\nAnother problem with mechanical linkage sensors is that the electronic components thereof are typically exposed to harsh environmental conditions (e.g., temperature extremes, water, dirt, salt) normally experienced by a vehicle traveling along a roadway. As a result of such exposure, the electronic components of the sensors can become corroded and fail to function properly. Due to one or both of these or other problems, one or more of the mechanical linkage sensors may be non-operational at any given time. Thus, regular inspection and replacement of such sensors is typically required.\nStill another disadvantage of mechanical linkage sensors is that the same are mounted separately from the other suspension components. As a result, additional time and effort is typically spent installing these components during the assembly process. Furthermore, additional effort is typically involved in creating a clearance area for mounting and operation of the mechanical linkage. Thus, such sensors disadvantageously require a significant amount of effort and space for mounting and operation.\nAs an alternative to mechanical linkage sensors, non-contact sensors that utilize sound or pressure waves traveling through a fluid medium, typically at an ultrasonic frequency, have been used in determining the relative position of one structural member to another. One example of such an application includes an ultrasonic sensor being used to determine a height of a fluid suspension member, such as an air spring. In such a use, the ultrasonic sensor is supported on one end member of the air spring and sends ultrasonic waves through the spring chamber of the air spring toward the opposing end member. The waves are reflected back by a suitable feature of the opposing end member and the distance therebetween is determined in a conventional manner.\nOne advantage of such an arrangement over mechanical linkages is that the ultrasonic sensor is at least partially sheltered from impacts and exposure. However, numerous disadvantages also exist with the use of ultrasonic sensors. One such disadvantage is that such sensors are relatively expensive which tends to undesirably increase production costs. Also, the replacement cost of a sensor that does get damaged by an impact or from exposure is likewise increased.\nAnother disadvantage is that ultrasonic sensors require a target that is suitable to reflect the ultrasonic waves back to the sensor for determining the distance therebetween. If such a target is not provided, the ultrasonic waves will not be reflected back properly and, thus, a correct determination of distance will not be possible. Thus, a target area must be provided for the proper operation of ultrasonic sensors. This can be particularly problematic, however, where the design constraints of a product limit the possibilities for including a target area. This is also a problem for existing products are being outfitted with ultrasonic sensors, where the existing products do not have a suitable target area."}
{"text": "Physical input devices (e.g., keyboard, mouse, stylus) are commonly used by users to perform actions (e.g., keystrokes, mouse movements). These actions may then be translated into data that can be provided to a computing device to represent human interactions with the computing device. Today, it is common for a user to have a multitude of computing devices that they use on a regular basis, and even simultaneously. For example, a single user may have a desktop computer, a laptop computer, a smart phone, and a tablet computer.\nAs a result, users may often regularly interface with more than one of these computing devices. For example, a user may be using a laptop computer for business purposes and a tablet computer for entertainment purposes. It would be cumbersome for a user to carry separate input devices for each computing device and physically move from one input device connected to one computing device to second input device connected to a second host computing device in order to utilize both computing devices. Some alternatives allow reusing a single input device for data entry to multiple computing devices. However, these solutions are not without their drawbacks.\nOne significant drawback with prior solutions is how the process of switching multiple data input devices from a first host computing device to a second host computing device is done. In some prior solutions, in order to switch both the keyboard device and the mouse device to a second host computing device, the user would have to perform two separate actions, which can be detrimental to the user experience. For example, in prior solutions, both the keyboard device and the mouse device are independently paired with each of the two or more host computing devices. When the user wants to switch from using a first host computing device to a second host computing device, the user would have to change the pairing for each of the keyboard device and the mouse device, individually.\nBased on the foregoing, there is a need in the art for improved methods and systems to more efficiently switch connections pairings input devices with multiple computing devices."}
{"text": "Subscriber devices are used as parts of wireless communication systems. Inherent in the operation of time division duplex (TDD) devices is a fundamental \"turn-around\" time which is the amount of time it takes for the devices to switch from a receive to a transmit mode and vice versa. Because of the turn-around time an amount of inefficiency is introduced into the speed of operation of the communication system. Base stations and infrastructure equipment for the wireless service are typically much more complex than the subscriber devices; therefore, the subscriber devices normally have larger turn-around times than the infrastructure equipment and the operation of the system is limited by the slowest subscriber device. More particularly, a guard time is set in the communication system in order to support the subscriber device with the slowest turn-around time. The speed of operation of the communication system is dependent on and determined by the turn-around time of the slowest subscriber device in the communication system. Therefore, the capacity of information throughput of the communication system is greatly reduced.\nThus, there is a need for a wireless communication system that provides a more efficient and greater information throughput by taking into account the turn-around time of the subscriber devices of the wireless communication system."}
{"text": "For critical infrastructure facilities, mitigation techniques for insider threats are primarily non-technical in nature and rely heavily on policies/procedures. Traditional access control measures (access cards, biometrics, PIN numbers, etc.) are built on a philosophy of trust that enables those with appropriate permissions to access facilities without additional monitoring or restrictions.\nMany authentication systems, such as those used in online banking and other web applications, operate on the basis of a virtual ID (e.g., session cookie), that is created after an initial authentication, typically via user name and a password. The security application essentially transforms a computing device into both a device that is associated with the user, and a location of the user. The session cookie itself becomes the item that is tied to the network source address of the computer, which becomes akin to a location. As GPS devices and smart phones continue to become smaller, less expensive and more powerful, the session cookies may also serve a dual purpose to identify the user location. Additionally, after initial access is granted, the authentication system may prompt the user to periodically re-enter access identification, particularly when changing applications, thereby disrupting the user.\nThus, there is a need for an authentication system with reduced disruption to the user that provides a high assurance of strong authentication.\nIntended advantages of the disclosed systems and/or methods satisfy one or more of these needs or provide other advantageous features. Other features and advantages will be made apparent from the present specification. The teachings disclosed extend to those embodiments that fall within the scope of the claims, regardless of whether they accomplish one or more of the aforementioned needs."}
{"text": "In biosciences and related fields, it can be useful to assay biological activity of micro-objects such as cells. Some embodiments of the present invention include apparatuses and processes for assaying biological activity in holding pens of a micro-fluidic device."}
{"text": "1. Field of the Invention\nThe present invention relates to memory storage devices, and, in particular, to buffering of data transfers for direct access block devices, such as solid state disks (SSDs).\n2. Description of the Related Art\nFlash memory is a type of non-volatile memory that is electrically erasable and re-programmable. Flash memory is primarily used in memory cards and USB flash drives for general storage and transfer of data between computers and other digital products. Flash memory is a specific type of electrically erasable programmable read-only memory (EEPROM) that is programmed and erased in large blocks. One commonly employed type of flash memory technology is NAND flash memory. NAND flash memory forms the core of the flash memory available today, especially for removable universal serial bus (USB) storage devices known as USB flash drives, as well as most memory cards. NAND flash memory exhibits fast erase and write times, requires small chip area per cell, and has high endurance. However, the I/O interface of NAND flash memory does not provide full address and data bus capability and, thus, generally does not allow random access to memory locations.\nThere are three basic operations for NAND devices: read, write and erase. The read and write operations are performed on a page by page basis. Page sizes are generally 2N bytes, where N is an integer, with typical page sizes of, for example, 2,048 bytes (2 kb), 4,096 bytes (4 kb), 8,192 bytes (8 kb) or more per page. Pages are typically arranged in blocks, and an erase operation is performed on a block by block basis. Typical block sizes are, for example, 64 or 128 pages per block. Pages must be written sequentially, usually from a low address to a high address. Lower addresses cannot be rewritten until the block is erased.\nA hard disk is addressed linearly by logical block address (LBA). A hard disk write operation provides new data to be written to a given LBA. Old data is over-written by new data at the same physical LBA. NAND flash memories are accessed analogously to block devices, such as hard disks. NAND devices address memory linearly by page number. However, each page might generally be written only once since a NAND device requires that a block of data be erased before new data is written to the block. Thus, for a NAND device to write new data to a given LBA, the new data is written to an erased page that is a different physical page than the page previously used for that LBA. Therefore, NAND devices require device driver software, or a separate controller chip with firmware, to maintain a record of mappings of each LBA to the current page number where its data is stored. This record mapping is typically managed by a flash translation layer (FTL) in software that might generate a logical-to-physical translation table. The flash translation layer corresponds to the media layer of software and/or firmware controlling an HDD.\nAssociated with each page is a spare area (typically 100-500 bytes) generally used for storage of error correction code (ECC) information and for storage of metadata used for memory management. The ECC is generally needed for detecting and correcting errors in the user data stored in the page, and the metadata is used for mapping logical addresses to and from physical addresses. As such, the additional bytes of memory are “hidden” from the user and are not available for storing user data. The first block (block 0) of a flash die is generally provided from the manufacturer error-free, and is commonly used by designers to include program code and associated metadata for block management.\nFor consumer applications, HDDs generally have data sectors that are sized in powers of two (e.g. 512 (29) bytes per sector). Flash memories structured with page sizes that are a multiple of the HDD sector size might efficiently work with the HDD system by storing multiple entire sectors in a page (e.g. a 4096 byte page can store eight 512 byte sectors). However, enterprise-based HDD systems generally do not use sectors sized by powers of two, but use larger sectors, generally either 520 or 528 bytes per sector instead of 512 bytes. Thus, typical flash memories perform inefficiently for enterprise applications since there are unused bytes in each page.\nTypically, for high capacity solid state disks (SSDs), several design tradeoffs might be considered when implementing a method to maintain a logical-to-physical translation table. These tradeoffs typically include: efficient random access memory (RAM) usage; efficient flash usage; fast address lookup for both read operations and write operations; fast write performance; and fast reconstruction of the translation table on device startup.\nSeveral techniques are known in the art for maintaining the logical-to-physical translation table. One such approach is known as direct page mapping, an example of which is described in the paper by Andrew Birrell & Michael Isard, et al., A DESIGN FOR HIGH-PERFORMANCE FLASH DISKS, ACM SIGOPS Operating Systems Review, Vol. 41, Issue 2, pp. 88-93, (April 2007), which is incorporated herein by reference in its entirety (hereinafter “Birrell”). Direct page mapping maintains a lookup table in RAM having an entry for each flash page, and a summary page for metadata at the end of each block, from which the logical-to-physical translation table may be reconstructed at startup. For example, a direct page mapped translation table might contain, for every LBA, a logical sector number corresponding to a physical block number and a physical page number. Thus, direct page mapping comprises a single-level logical-to-physical translation. The summary page for each block might contain the LBA and valid bits for each page in the block so that the translation table can be reconstructed at startup. Thus, the direct page mapping scheme requires a large amount of RAM (on the order of 1-2 MB per GB of user storage) to store the translation table, which can become burdensome for higher capacity SSDs.\nAnother approach is known as block mapping. Block mapping generally classifies blocks as either data blocks (D-blocks) or update blocks (U-blocks). The total size of the D-blocks is the effective storage space for user data while U-blocks are invisible to users. Generally, when a write command cannot be accommodated in the D-block corresponding to the LBA, a U-block is allocated to receive the new data and the old data in the D-block is invalidated. Subsequent writes to that D-block will be received by the allocated U-block. When the U-block becomes full, another U-block might be allocated, or the U-block might be merged with the original D-block. Thus, block mapping maintains a lookup table in RAM that maps a logical block to a physical block. Block mapping lacks a page-level map, instead relying on the typical case that data is stored in sequential order within the block. For example, a block mapped translation table might contain a logical sector number corresponding to a logical block number and a logical page number. The logical block number can be translated into a physical block number and the logical page number might correspond to a physical offset within the physical block. Thus, block mapping comprises a two-level logical-to-physical translation. The size of the translation table is proportional to the number of blocks in the flash memory, thus requiring less RAM than a page mapped translation table.\nHowever, because block mapping does not have a page-level map, the flash media may be inefficiently utilized when the data access workload is non-sequential. For non-sequential data access workloads, block mapping might require data to be copied and re-written numerous times to maintain the correct mapping. An example of block mapping is described in the paper by Jeong-Uk Kang & Heeseung Jo, et al., A SUPERBLOCK-BASED FLASH TRANSLATION LAYER FOR NAND FLASH MEMORY, Proceedings of the 6th ACM & IEEE International Conference On Embedded Software, pp. 161-170, (Oct. 22-25, 2006), which is incorporated herein by reference in its entirety (hereinafter “Kang”).\nA third approach for maintaining the logical-to-physical translation table is known as a superblock mapping scheme. Superblock mapping groups together a set number of adjacent logical blocks into a Superblock. Superblock mapping maintains a page global directory (PGD) in RAM for each Superblock. Page middle directories (PMDs) and page tables (PTs) are maintained in the spare areas of the flash pages. Each LBA can be divided into a logical block number and a logical page number, with the logical block number comprising a superblock number and a PGD index offset. The logical page number comprises a PMD index offset and a PT index offset. Each entry of the PGD points to a corresponding PMD. Each entry of the PMD points to a corresponding PT. The PT contains the physical block number and the physical page number of the data. To translate a logical address to a physical address in Superblock mapping, a module must access RAM to read the PGD, access flash to read the PMD, access flash to read the PT, and access flash to access the requested data address. Super-block mapping, thus, comprises a four-level logical-to-physical translation and provides page-mapping.\nThe PMD's and PT's are stored in the spare areas of the flash pages to provide page-mapping without using an excessive amount of RAM. However, because the spare area is used to store page-level mapping information, less memory is available for error correction codes (ECC). Further, the limited amount of memory available in the spare area precludes storing complicating mapping information. Finally, reconstruction of the translation table at startup can be time-intensive. An example of a superblock mapping scheme is described in Kang.\nAs described previously, for write operations, NAND devices store the new data for the LBA on a new page, unlike hard disk drives (HDDs) that can rewrite individual physical sectors. Thus, a NAND device generally requires that a block be erased before new data can be written to the block. Further, as described above, often a NAND device will write new data for a given LBA to an erased page that is a different physical page from the page previously used for that LBA. Thus, NAND devices also generally require the device driver software or the separate controller chip periodically initiate a process to erase data that is “stale” or out-of-date. As would be apparent to one of skill in the art, without periodically erasing out-of-date data, the flash memory would fill up with data that is mostly out-of-date. This inefficiency would reduce the realized flash memory capacity because less current data could be stored. Therefore, device driver software or controller chips generally periodically run a “garbage collection” routine adapted to provide efficient flash memory utilization by erasing out-of-date blocks. An example of a garbage collection routine is described in Kang. Garbage collection routines impact performance of the flash memory system by utilizing processor resources and potentially delaying write operations to the flash media.\nHowever, NAND device blocks can be erased relatively few times before device failure (typically on the order of 100,000 erasures). Therefore, over the operational life of an SSD, blocks of flash memory will fail and become unusable. Thus, the device driver software or the separate controller chip should minimize the number of erasures, and must also maintain a record of bad blocks. For example, device driver software or controller chips might implement wear leveling to spread the erasing and writing of blocks over the entire flash memory to avoid repeatedly erasing and writing a given subset of blocks."}
{"text": "The field of the invention is handlebar grips and the invention relates more particularly to handlebar grips which are held to the handlebar in such a way that they do not turn during use. Such handlebar grips are used on bicycles, all-terrain vehicles (ATVs), watercraft and snowmobiles. Handlebar grips, particularly for use on mountain bikes, are preferably soft on the outside so that they can be held without irritating the rider's hands over rough terrain. However, soft grips are often difficult to adhere to the handlebars without turning and it is advantageous that the handlebar grips do not turn with respect to the handlebar.\nVarious designs have been patented in an attempt to provide a securely held soft grip. One such design is shown in the Kuipers, et al. U.S. Pat. No. 5,280,735. This patent shows a two-piece slip-resistant grip which has an outer nylon lattice housing member with two end clamps. This lattice and end clamp assembly is placed over a softer inner sleeve which extends through the openings in the lattice and is clamped in place by the end cap and flanged end. This assembly is limited to a design where the hand contacting portion is in the shape of diamonds which pass through the lattice. Also, this grip results in the soft inner portion contacting the handlebar and with extended time the soft portion tends to form a bond with the handlebar and is very difficult to remove.\nU.S. Pat. No. 3,344,684 to Steere, Jr., et al., shows a handlebar grip which has an inner portion with a transparent or translucent outer portion adhered thereto. The inner portion is slipped on the handlebar without any tightening means. Another two-part grip is shown in U.S. Pat. No. 4,416,166 to Jannard, et al. which has a two-part grip having an outer sleeve fitted over an inner sleeve. These two parts are interlocked by their matching designs.\nU.S. Pat. No. 4,535,649 shows a foam surfaced handlebar grip which has a foam sleeve which is held over a hard polymer tube by an end cap and a coupling member. A cork handlebar grip is shown in U.S. Pat. No. 605,626 to Blanchard, Jr. This patent utilizes a slotted tapered pipe which has a nut at each end which captures the cork grip portion."}
{"text": "Biopsy samples have been obtained in a variety of ways using various devices. An exemplary biopsy device is the MAMMOTOME device from Ethicon Endo-Surgery, Inc. of Cincinnati, Ohio. Further exemplary biopsy devices are disclosed in U.S. Pat. No. 5,526,822, entitled “Method and Apparatus for Automated Biopsy and Collection of Soft Tissue,” issued Jun. 18, 1996; U.S. Pat. No. 6,086,544, entitled “Control Apparatus for an Automated Surgical Biopsy Device,” issued Jul. 11, 2000; U.S. Pub. No. 2003/0109803, entitled “MRI Compatible Surgical Biopsy Device,” published Jun. 12, 2003; U.S. Pub. No. 2007/0118048, entitled “Remote Thumbwheel for a Surgical Biopsy Device,” published May 24, 2007; U.S. Provisional Patent Application Ser. No. 60/869,736, entitled “Biopsy System,” filed Dec. 13, 2006; U.S. Provisional Patent Application Ser. No. 60/874,792, entitled “Biopsy Sample Storage,” filed Dec. 13, 2006; and U.S. Non-Provisional patent application Ser. No. 11/942,785, entitled “Revolving Tissue Sample Holder for Biopsy Device,” filed Nov. 21, 2007. The disclosure of each of the above-cited U.S. patents, U.S. Patent Application Publications, U.S. Provisional Patent Applications, and U.S. Non-Provisional patent application is incorporated by reference herein. While many of the foregoing biopsy devices are configured to obtain biopsy samples from breast tissue, biopsy samples may also be obtained from various other locations.\nVarious biopsy devices may be designed to work with X-ray, ultrasound, and magnetic resonance imaging (MRI) as imaging modalities. For instance, various components for interfacing biopsy devices with various imaging systems are disclosed in the following: U.S. Pub. No. 2005/0261581, entitled “MRI Biopsy Device,” published Nov. 24, 2005; U.S. Pub. No. 2005/0277829, entitled “MRI Biopsy Apparatus Incorporating a Sleeve and a Multi-Function Obturator,” published Dec. 15, 2005; U.S. Pub. No. 2005/0283069, entitled “MRI Biopsy Device Localization Fixture,” published Dec. 22, 2005; U.S. Pub. No. 2007/0167736, entitled “MRI Biopsy Apparatus Incorporating an Imageable Penetrating Portion,” published Jul. 19, 2007; U.S. Pub. No. 2006/0241385, entitled “Guided Disposable Fiducial for Breast Biopsy Localization Fixture,” published Oct. 26, 2006; U.S. Pub. No. 2006/0258956, entitled “MRI Biopsy Device,” published Nov. 16, 2006; U.S. Pub. No. 2007/0255168, entitled “Grid and Rotatable Cube Guide Localization Fixture for Biopsy Device,” published Nov. 2, 2007; and U.S. Pub. No. 2007/0255170, entitled “Biopsy Cannula Adjustable Depth Stop,” published Nov. 1, 2007; and US Pub No. 2008/0015429, “MRI Biopsy Device” published Jan. 17, 2008. The disclosure of each of the foregoing published patent applications is incorporated by reference herein.\nIt may be desirable in some settings to use one or more imaging modalities other than X-ray, ultrasound, or MRI before, during, or after a biopsy procedure. For instance, an alternative imaging modality may include positron emission tomography (PET) imaging. In a mammography context, such imaging may be referred to as positron emission mammography (PEM). Instead of scanning the entire body, PEM may be used as a special form of PET for imaging breasts and other small body parts. This may allow for a more detailed image of abnormal tissue. In a PEM context, the patient may be injected with an intravenous substance called FDG (fluorodeoxyglucose), which is a glucose analog, which may accumulate in glucose avid cells. This substance may carry a positron emitting radioactive isotope. One or more detectors may be used to capture emission of positrons emitted by such an isotope (e.g., by capturing resulting gamma photons) to ultimately produce an image. Alternatively, any other substances may be injected into a patient, as a tracing agent for PEM imaging or otherwise. An exemplary PEM system may include the PEM FLEX SOLO II system by Naviscan PET Systems, Inc. of San Diego, Calif.\nAnother alternative imaging modality may include breast-specific gamma imaging (BSGI). In a use of BSGI, a patient may be injected with a radiotracer (e.g., Technicium isotope T-99), and a BSGI camera may be used to capture gamma radiation emitted by such a tracer. Cancerous cells may have a higher tendency to absorb certain gamma emitting radiotracers, which may result in cancerous lesions standing out under BSGI imaging. BSGI imaging may thus provide distinction between cancerous tissue and non-cancerous tissue based on cellular activity rather than being based on tissue density. An exemplary BSGI system may include the DILON 6800 by Dilon Technologies of Newport News, Va.\nVarious biopsy site marker devices are disclosed for use in marking biopsy sites. One or more marker devices are disclosed in U.S. Pub. No. 2005/0228311, entitled “Marker Device and Method of Deploying a Cavity Marker Using a Surgical Biopsy Device,” published Oct. 13, 2005; U.S. Pat. No. 6,996,433, entitled “Imageable Biopsy Site Marker,” issued Feb. 7, 2006; U.S. Pat. No. 6,993,375, entitled “Tissue Site Markers for In Vivo Imaging,” issued Jan. 31, 2006; U.S. Pat. No. 7,047,063, entitled “Tissue Site Markers for In Vivo Imaging,” issued May 16, 2006; U.S. Pat. No. 7,229,417, entitled “Methods for Marking a Biopsy Site,” issued Jun. 12, 2007; U.S. Pat. No. 7,044,957, entitled “Devices for Defining and Marking Tissue,” issued May 16, 2006; U.S. Pat. No. 6,228,055, entitled “Devices for Marking and Defining Particular Locations in Body Tissue,” issued May 8, 2001; and U.S. Pat. No. 6,371,904, entitled “Subcutaneous Cavity Marking Device and Method,” issued Apr. 16, 2002. The disclosure of each of the above-cited U.S. patents and U.S. Patent Application Publications is incorporated by reference herein."}
{"text": "1. Field of the Invention\nThe present application relates to magnetic resonance (MR) imaging system, and more specifically to MR imaging system which allows interactive selection of imaging planes during imaging.\n2. Description of Related Art\nIn conventional magnetic resonance (MR) imaging systems, the three dimensional location of a desired region of the patient to be imaged must be provided to the MR imaging system before obtaining the image. In addition to the location, the three dimensional orientation must also be provided. Together these define an `imaging plane`. These locations and orientations typically are either calculated manually, or prescribed graphically, by drawing a line on a previously acquired image. Even though the computations may not be difficult, it makes it cumbersome to produce several images at different orientations. If images are being rapidly acquired in a continuous imaging mode, it is desirable that the operator keep his eyes on the new images as they are displayed, and these imaging-plane prescription techniques do not allow this. Also, the graphic prescription technique described above will only define imaging planes orthogonal to the previously acquired image.\nWhile acquiring images of the heart it is often desirable to adjust the imaging plane interactively while searching for anatomy such as the coronary arteries. Also, with new open magnet MR Imaging systems, as described in U.S. Pat. No. 5,365,927 issued Nov. 22, 1994 \"Magnetic Resonance Imaging System With Pointing Device\" Roemer et al., assigned to the present Assignee, and hereby incorporated by reference; it is possible to perform medical procedures on a subject and simultaneously acquire a series of continuous MR images. When a medical procedure, such as surgery, is being performed, it is very desirable to minimize the time required to calculate location and orientation of `imaging planes` to acquire images.\nCurrently there is a need for an MR imaging system which provides MR images along selected orientations through selected internal structures in the heart of a patient, in which an operator may easily indicate location and orientation of an imaging plane, and acquire an image of the subject from that plane."}
{"text": "Large format inkjet printers use vacuum tables to hold down foamboard, cardboard and other inflexible or semi-flexible print media for printing. High capacity vacuum pumps are used to develop the hold down forces needed to keep large sheets of such media flat during printing.\nThe same part numbers designate the same or similar parts throughout the figures. The figures are not necessarily to scale. The size of some parts may be exaggerated to better illustrate the example shown."}
{"text": "1. Field of the Invention\nThis invention relates to a percussive tool, and in particular but not exclusively to a percussive tool for use in drilling operations utilizing drilling fluid or xe2x80x9cmudxe2x80x9d.\n2. Description of the Prior Art\nIt has long been appreciated that providing a percussive action to a drill bit can increase the drilling rate, particularly when drilling through hard rock. Air driven percussion drills have been used successfully for a number of years in the mining industry. Also, hydraulic fluid driven percussion tools are well established in the construction industry. However, attempts to provide a percussion drill utilizing drilling fluid or xe2x80x9cmudxe2x80x9d as the working fluid have met with difficulties; the evasive nature and high solids content of drilling fluid tends to damage the valve mechanisms necessary to provide the percussive effect.\nIt is among the objectives of embodiments of the present invention to provide a percussion drill which will operate successfully and reliably utilizing drilling fluid or mud as the working fluid.\nAccording to the present invention there is provided a percussion drill comprising:\na fluid transmitting body;\na drill bit support;\na mass movable relative to the body for impacting on the drill bit support;\nmeans associated with the mass for creating a fluid pressure force on said mass; and\na rotating valve for controlling flow of fluid through the body to produce a varying fluid pressure force on the mass and induce acceleration of the mass.\nThe use of a rotating valve facilitates use of the drill in applications where the working fluid contains solids, such as drilling fluid.\nAccording to another aspect of the present invention there is provided a percussive drilling method comprising the steps:\nproviding a drill having a fluid transmitting body, a drill bit support, a mass movable relative to the body for impacting on the drill bit support and a rotating valve for controlling flow of fluid through the body; and\npassing drilling fluid through the body to the drill bit support via the rotating valve to vary the flow of the fluid through the body and produce a varying fluid pressure force on the mass to induce acceleration of the mass.\nPreferably, the valve rotates around a longitudinal axis, and most conveniently around the central longitudinal axis of the body. In other embodiments the valve may rotate around a transverse or lateral axis.\nThe valve may be provided separately of the means for creating fluid pressure force on the mass, or may be integral therewith.\nPreferably also, the valve comprises two portions, each defining a fluid port, such that relative rotation of the parts varies the alignment of the ports and varies the flow area defined thereby. Most preferably, one portion is rotatable relative to the body. The valve ports may be in the form of slots on a common axis. In an alternative embodiment, one of the valve portions provides fluid communication with alternative fluid paths through the body, one of said fluid paths providing fluid communication with the means for creating fluid pressure force on the mass, and another of said paths by-passing said means.\nPreferably also, the valve permits fluid flow through the body in all valve configurations, to assure a continuing supply of drilling fluid to the drill bit.\nPreferably also, the mass is spring mounted in the body, to provide a return force; the fluid pressure force will tend to induce acceleration of the mass in one direction, and the return action of the spring will accelerate the mass in the opposite direction, when the spring force is greater than the fluid pressure force. In other embodiments, the return force may simply be gravity, or may be an opposing fluid pressure force.\nThe mass may define a flow passage therethrough, and the flow passage may define a restriction such that fluid flowing through the mass is subject to a pressure drop, creating a pressure force across the restriction.\nThe drill bit support will typically be adapted for mounting a drill bit to the body and the mass will impact directly on the support. However, in some embodiments the mass may not impact directly on the support; there may be an intermediate member or other force transmission means therebetween.\nPreferably also, the drill bit support is spring mounted in the body, and is preferably splinted to the body.\nPreferably also, the body is adapted for mounting to a drill string.\nPreferably also, the drill includes means for driving the valve, such as a valve motor, and most preferably a drilling fluid driven positive displacement motor.\nAccording to a further aspect of the present invention there is provided a percussion tool comprising:\na fluid transmitting body;\na bit support;\na mass movable relative to the body for impacting on the bit support;\nmeans associated with the mass for creating a fluid pressure force on said mass; and\na valve for providing a continuing but varying flow of fluid through the body to produce a varying fluid pressure force on the mass and induce acceleration of the mass.\nThe use of a valve which is not required to stop flow through the body, that is the valve is not required to seal, facilitates use of the drill in applications where the working fluid contains abrasives or solids, such as drilling fluid; as the valve is not required to be fluid tight, clearances, materials and other aspects of the valve specification may be selected with the primary aim of withstanding operation using abrasive working fluid.\nAccording to a still further aspect of the present invention there is provided a method of providing a percussive action in a downhole tool, the method comprising the steps:\nproviding a tool having a fluid transmitting body, a bit support, a mass movable relative to the body for impacting on the bit support and a valve for controlling flow of fluid through the body; and\npassing drilling fluid through the body and the valve to provide a varying but continuing flow of fluid through the body and producing a varying fluid pressure force on the mass to induce acceleration of the mass.\nThe bit support may be integral with a bit, or may provide mounting or support for a separable bit. The bit may be a drill bit, a chisel or a hammer."}
{"text": "1. Field of the Invention\nThe present invention concerns an inductive proximity sensor of the type including a coil for sensing metal objects forming part of an oscillatory circuit, an oscillator driving the oscillatory circuit at a constant frequency, a circuit for processing the signal generated by the oscillatory circuit, an output circuit and a device for tuning the frequency of the oscillatory circuit to the frequency of the oscillator.\n2. Description of the Prior Art\nThe signal generated by oscillatory circuit varies according to the distance of a metal object. In a sensor of this kind it is desirable to tune the frequency of the oscillatory circuit to that of the oscillator, to improve the sensitivity of the sensor. If the natural frequency of the oscillatory circuit is varied by varying the capacitance in this circuit the range of adjustment is limited by the fact that the components currently available--variable capacitors and variable capacitance semiconductor diodes--currently offer only a small range of variation, in the order of a few hundred pF. The frequency change that can be obtained is therefore very small, unless extremely large capacitors are used.\nA frequency divider can be inserted between the oscillator and the oscillatory circuit. In this case, however, the resulting flexibility is inversely proportional to the required accuracy of adjustment. Also, to obtain the required range of adjustment a very high frequency must be used, and therefore very fast components; the resulting implementation cost is therefore high.\nAn object of the invention is to remedy these drawbacks by employing a simple solution to obtain a wide range of adjustment, well suited to inductive proximity sensors."}
{"text": "1. Field of the Invention\nThe present invention relates generally to earth boring rotary rock bits and, more particularly, to rock bits having one piece bit bodies.\n2. Description of the Prior Art\nA rotary rock bit usually is manufactured from a plurality of subassemblies. Each subassembly comprises mainly a leg having a journal pin formed on the lower extremity thereof. A cone type cutter is rotatively supported on the journal pin and is usually retained thereon by a plurality of ball bearings. Fluid nozzles and self-contained lubrication systems are also contained within each leg. The leg subassemblies are then welded together to form a bit body.\nA serious shortcoming with rock bits made in this manner is that a true geometry bit is practically impossible to manufacture. Because of the tolerance buildup in each of the subassemblies and because of the final welding process of the three legs, unwanted bit skew is incorporated into the final product. Moreover, the cone heights of the assembled products are not always the same. In addition, the journal offset and angle are not accurately held using the conventional assembly techniques.\nAnother problem with manufacturing rock bits from leg subassemblies is that many additional steps are required in the process. First of all, the leg seams must be machined with a 120.degree. milling machine, and dowel holes must be drilled. Moreover, the journal pin bearing surfaces must be machined and having the journal pins integral with the legs makes the machining operation lengthy and cumbersome. In addition, the bearing surfaces of each journal pin must necessarily be carburized while the rest of the leg surface be painted over to prevent the remaining surface from being curburized. However, carbon leakage does occur on this surface and this leakage is deleterious to the subsequent operations of the legs.\nBecause of the problems associated with the above-mentioned leg subassemblies, one piece bit bodies have been proposed. Also, separate journal pins have been proposed in order to carburize and machine the pins prior to assembling them to the legs.\nU.S. Pat. No. 4,043,411 shows one method of having a separate journal pin secured to a leg. However, the shortcoming found in the method of manufacture of the rock bit legs of the referenced patent is that the welding is accomplished on the interior side of the leg, i.e., the cutter receiving surface side of the leg. Such a method of manufacture severely restricts the design capabilities of the journal pin. Moreover, with such a restriction, the cone cannot be preassembled with the journal pin prior to welding the pin to the leg and may require secondary operations to clean welded areas prior to final assembly.\nAnother U.S. Pat. No. 4,098,150 also discloses separate journal pins. However, the structure disclosed therein does not lend itself to a one piece bit body construction.\nU.S. Pat. No. 3,850,256 discloses a modified one bit body construction; however, very few of the problems discussed above are solved with the structure disclosed therein."}
{"text": "The electronic game industry has seen a dramatic evolution from the first electronic ping-pong game (\"pong\") to the state of modern games and consumer home electronics. In general, hardware advances that have increased processing power and reduced cost has fueled this evolution. The increased availability of low cost processing power, as well as consumer expectation for improved game content, demands that new games be developed to take advantage of this processing power. This can be seen especially in the new 64 bit processing devices such as the Nintendo 64.TM. and the processing power available in home personal computer games and/or in arcade game platforms. These new hardware platforms are so powerful that a whole new genre of games has to be developed in order to fully utilize the hardware.\nElectronic game input, traditionally, has been limited to joy sticks, button paddles, multi-button inputs, trackballs and even a gyro mouse that has a gyroscope means for determining the orientation of the mouse. Recently Nintendo has deployed a \"rumble\" device to provide vibratory feedback to game console users. Traditional computer input means are well know to those in the arts and require no further discussion. The gyro-mouse, in the context of the present invention, however, deserves some further discussion.\nThe gyro-mouse, provided in U.S. Pat. No. 5,138,154 to Hotelling et al., the relevant portions herein incorporated by reference in their entirety, provides a means for using the gyroscopic effect in a computer input device to recover user input. The gyro-mouse provides a gyroscope contained within a ball so that ball may be rotated. This rotation translates into two dimensional or three dimensional motion for software receiving the gyro-mouse input to display on a computer screen. Thus, the gyro-mouse is somewhat an extension of the track ball paradigm for a computer input device.\nThe gyroscopic effect has also been harnessed for practical commercial applications. One of the more interesting gyroscopic effects is brought about through the principal of conservation of angular momentum. As witnessed in gyroscopic phenomena, a gyroscope creates a force at right angles to a force that attempts to \"topple\" the gyroscope. Thus, a gyroscope when left alone or mounted in a double gimbal arrangement allowing the gyroscope to move freely in both axes, will resist movement and/or attempt to hold its own angular position. Gyroscopes are also known to have precession due to the earth's effect on the gyroscope. Gyroscope procession is not especially pertinent to the present invention, however, its principals and mathematical proofs and formula are herein incorporated by reference.\nThe navigational arts also provide a means for harnessing gyroscopic phenomena to determine the inertial position of a vehicle such as an aircraft. In an inertial navigation system, the gyroscope is mounted in a double gimballed arrangement and allowed to rotate without resistance in all directions. As the aircraft turns, rotates, and/or changes direction the gyroscopic effect keeps the inertial navigation gyroscope at the same angle. High precision means are used to determine how much the gyrostat has rotated, in actuality the aircraft rotating around the gyroscope, and this measurement in combination with high precision accelerometers provides a means for tracking the change in an aircraft direction. This instrumentality in conjunction with precision timing and velocity measurements provides a means for continuously determining an aircraft navigational position.\nIn another application of the gyroscopic effect, a large gyroscope can be used to create an effect that in some aspects is the reverse that of an inertial navigation system. Here, a large gyrostat mass (the flywheel) can be use to stabilize or position certain objects such as spacecraft. In the spacecraft application, such as in U.S. Pat. No. 5,437,420, the relevant portions herein incorporated by reference in their entirety, large flywheels and high torque motors and brakes are used to topple the flywheel. The spacecraft then feels a moment of thrust at right angles to the torque that is applied to the gyrostat. This way, and in others such as the \"pure\" inertia of rotation of causing a flywheel mass to accelerate or decelerate rotation, spacecraft attitude may be changed through gyrostatic means. In other stabilization applications, gyrostats are used to stabilize platforms such as cameras and other precision instruments, in general by attaching a gyrostat to the instrument platform.\nGyrostats have been used in conjunction with wheels to provide linear propulsion. Through a systems of gears and linkages, U.S. Pat. No. 5,090,260, incorporated herein by reference in its entirety, provides a means for translating the gyrostatic toppling effect into a linear force for propulsion."}
{"text": "1. Field of the Invention\nThe present invention relates to an imaging lens for use in imaging equipment using an imaging element such as CCD (Charge Coupled Device) and CMOS (Complementary Metal Oxide Semiconductor).\n2. Description of Related Art\nAn imaging element such as CCD and CMOS, recently, has been greatly reduced in size and increased in the number of pixels. For this reason, the main body of the imaging equipment and the lens mounted on the main body of the imaging equipment have been also required of reduced size and higher performance. Also, to correspond to the imaging element with pixels increased in number, the imaging lens is also requested to have a telecentric characteristic, that is, it is requested that an incident angle onto the imaging element is near to parallel to an optical axis (namely, the incident angle onto the imaging element is near to perpendicular to the imaging element). In the related art, for example, most of imaging lenses mounted on a cellular phone with a camera are composed of three lenses; however, in order to cope with the increased number of pixels, there is a tendency toward larger number of lenses. In an imaging lens disclosed in JP-2005-208236, a second lens disposed second from the object side is a negative lens, the surface of which on the object side is concave.\nHowever, in the imaging lens disclosed in JP-2005-208236, there is room for improvement with respect to the chromatic aberration thereof and the like. Also, in an imaging lens composed of four lenses, in order to secure the telecentric characteristic, it is expected that an aperture stop is disposed on the object side as much as possible. In this case, when, as in the imaging lens provided by JP-2005-208236, the object side surface of the second lens is concave, the manufacturing sensitivity thereof is high and, therefore, for example, when the position of the lens is shifted in the manufacturing process thereof, the image surface thereof is easy to vary."}
{"text": "1. Field of the Invention\nThis invention relates to improved liquid polysulfide polymer based adhesive compositions containing novel adhesive additives. More particularly, this invention relates to improved liquid polysulfide polymer based adhesive compositions containing novel adhesive additives prepared by modifying polysulfide and polyester polymers with abietic acid and polysulfide polymers with epoxy groups.\n2. Description of the Prior Art\nCurable, liquid polysulfide polymer based sealant and caulking compositions have long been known in the art. They are based on curable, liquid, mercaptan terminated polysulfide polymers. The structure and preparation of these polythiopolymercaptan polymers are disclosed in U.S. Pat. No. 2,466,963. These liquid, mercaptan terminated polymers are used extensively in sealant and caulking compositions, especially in the building, automotive, and marine industries; for castings, e.g., solid rocket propellants; for leather and textile impregnating agents; for adhesives, coatings, etc. Because of the wide range of applications in which these polysulfide polymer based sealing and caulking compositions are used, they must be capable of bonding to various substrates, e.g., metal, aluminum, glass, concrete, wood, etc. However, the adhesive qualities of polysulfide polymers are such as to usually require the use of an adhesive additive in the polysulfide polymer based sealant composition in order to insure adequate bonding of the cured sealant to the substrate.\nAdhesive additives of the prior art such as those of a phenolic nature provided only a polar, rather than a chemical, type bonding of the cured sealant to the substrate. Thus, sealant compositions containing phenolic type adhesive additives were found to be susceptible to attack by polar solvents such as water. The polar solvent tended to break down the polar bond thereby impairing the utility of the sealant composition. More recently, as disclosed in U.S. Pat. Nos. 3,297,473; 3,312,669, and 3,328,451, certain organo-silanes e.g., (mercapto alkyl amino alkyl alkoxyl silanes) have been proposed to overcome these problems. However, while the organo-silane adhesive additives have been found to protect the sealant composition bonds from attack by polar solvents, they do have certain disadvantages. Thus, the organo-silane adhesive additives are economically costly, somewhat deficient in storage stability, and exhibit erratic adhesion results with certain substrates, such as concrete.\nThe object of the present invention, therefore, is to provide a novel liquid polysulfide polymer based caulking or sealant composition which will allow for the attainment of a strong and lasting bonding of the sealant to the substrates to be treated therewith.\nAnother object of the present invention is to provide a novel liquid polysulfide polymer based caulking or sealant composition containing adhesive additives which are less costly, exhibit improved storage stability, and which exhibit more reproducible and uniform adhesion characteristics to a wider range of substrates.\nOther objects of the invention will become apparent to those skilled in the art from a consideration of the following detailed description."}
{"text": "This invention relates to motor vehicle manufacturing and more particularly to a high density welding machine to facilitate manufacture of subassembly components of motor vehicle bodies.\nIn the mass production of motor vehicles on an assembly line basis, it is imperative that the various subassembly components come together for the final assembly process in an efficient and orderly manner. The current subassembly procedures, while generally satisfactory, tend to be manpower-intensive, tend to require large amounts of factory floor space, sometimes present maintenance complications, require excessive capital expenditures, and may present safety and/or environmental concerns."}
{"text": "Most of any one person's DNA, some 99.9 percent, is exactly the same as any other person's DNA. The roughly 0.1% difference in the genome sequence accounts for a wide variety of the differences among people, such as eye color and blood group. Genetic variation also plays a role in whether a person is at risk for getting particular diseases or whether a person is likely to have a favorable or adverse response to a particular drug. Single gene differences in individuals have been associated with elevated risk for acquiring a variety of diseases, such as cystic fibrosis and sickle cell disease. More complex interrelationships among multiple genes and the environment are responsible for many traits like risk for some common diseases, such as diabetes, cancer, stroke, Alzheimer's disease, Parkinson's disease, depression, alcoholism, heart disease, arthritis and asthma.\nGenetic-based diagnostic tests are available for several highly penetrant diseases caused by single genes, such as cystic fibrosis. Such tests can be performed by probing for particular mutations or polymorphisms in the respective genes. Accordingly, risk for contracting a particular disease can be determined well before symptoms appear and, if desired, preventative measures can be taken. However, it is believed that the majority of diseases, including many common diseases such as diabetes, heart disease, cancers, and psychiatric disorders, are affected by multiple genes as well as environmental conditions. Thus, diagnosis of such diseases based on genetics is considerably more complex as the number of genes to be interrogated increases.\nRecently, through a variety of genotyping efforts, a large number of polymorphic DNA markers have been identified, many of which are believed to be associated with the probability of developing particular traits such as risk of acquiring known diseases. Exemplary polymorphic DNA markers that are available include single nucleotide polymorphisms (SNPs) which occur at an average frequency of more than 1 per kilobase in human genomic DNA. Many of these SNPs are likely to be therapeutically relevant genetic variants and/or involved in genetic predisposition to disease. However, current methods for genome-wide interrogation of SNPs and other markers are inefficient, thereby rendering the identification of useful diagnostic marker sets impractical.\nThe ability to simultaneously genotype large numbers of SNP markers across a DNA sample is becoming increasingly important for genetic linkage and association studies. A major limitation to whole genome association studies is the lack of a technology to perform highly-multiplexed SNP genotyping. The generation of the complete haplotype map of the human genome across major ethnic groups will provide the SNP content for whole genome association studies (estimated at about 200,000-300,000 SNPs). However, currently available genotyping methods are cumbersome and inefficient for scoring the large numbers of SNPs needed to generate a haplotype map.\nThus there is a need in the art for methods of simultaneously interrogating large numbers of gene loci on a whole genome scale. Such benefits will affect the genomic discovery process and the genetic analysis of diseases, as well as the genetic analysis of individuals. This invention satisfies this need and provides other advantages as well. This invention describes and demonstrates a method to perform large scale multiplexing reactions enabling a new era in genomics."}
{"text": "The present invention relates generally to machine vision, and more particularly to machine vision systems for forming a one dimensional digital representation of a low information content scene, e.g., a scene that is sparsely illuminated by an illumination plane, and the one dimensional digital representation is a projection formed with respect to columns of a rectangular pixel array.\nA well-known method for acquiring 3D range images includes steps of: providing a light source with line generating optics to illuminate a single plane of a scene, positioning a digital camera to view the light plane such that objects illuminated by the light source appear in the optical image formed by the camera lens, capturing a digital image of the scene, processing the digital image to extract the image coordinates of points in the scene illuminated by the light source, and processing the image coordinates according to the triangulation geometry of the optical system to form a set of physical coordinates suitable for measurement of objects in the scene.\nA major limitation associated with such a conventional machine vision process is that a 2-dimensional intensity image of substantial size must be captured by the digital camera for each and every line of physical coordinates formed by the system. This can make the time to capture the 3D image of a scene as much as 100 times longer than the time required to acquire an intensity image of the same size scene, thereby rendering laser-line based 3D image formation methods too slow for many industrial machine-vision applications."}
{"text": "The following references are also believed to represent the state of the art:\nUS Published Patent Application 2002/0087885 of Peled, et al.;\nUS Published Patent Application 2007/0220592 of Muehlbauer;\nU.S. Pat. No. 4,907,080 to Campbell, at al.;\nU.S. Pat. No. 5,870754 to Dimitrova, et al.;\nU.S. Pat. No. 6,466,670 to Tsuria, et al.;\nU.S. Pat. No. 6,633,651 to Hirzalla, et al.;\nU.S. Pat. No. 6,961,467 to Sirivara, et al.;\nU.S. Pat. No. 6,968,337 to Wold;\nPCT Published Patent Application WO2005/043359 of Macrovision\nCorporation;\nPCT Published Patent Publication WO2008/059393 of NDS Limited;\nPCT Published Patent Publication WO2008/096281 of NDS Limited;\nUK Published Patent Application GB2369203 of Assertion Limited;\nPaper entitled “Efficient video similarity measurements with video signature” by Sen-ching S. Cheung and Avideh Zakhor, University of California, Berkeley;\nPaper entitled “Video similarity detection with video signature clustering” by Sen-ching S. Cheung and Avideh Zakhor, University of California, Berkeley;\nPaper entitled “Fast similarity search and clustering of video sequences on the world-wide web” by Sen-ching S. Cheung, Center of Applied Scientific Computing Lawrence Livermore National Laboratory and Avideh Zakhor, University of California, Berkeley;\nPaper entitled “Robust video signature based on ordinal measure” by Xian-Sheng Hua, Xian Chen and Hong-Jiang Zhang, Microsoft Research, Asia;\nPaper entitled “Fast video matching with signature alignment” by Timothy C. Hoad and Justin Zobel, RMIT University, Melbourne, Australia;\nPaper entitled “A scalable signature scheme for video authentication” by Pradeep Atrey, Wei-Qi Yan, Mohan S. Kankanhalli, National University of Singapore;\nPress Release Mar. 13, 2006 of Motion DSP for Ikena Copyright;\nPaper entitled “Call for proposals on image & video signature tools” by International Organisation for Standardization, ISO/IEC JTC1/SC29/WG11, Coding Of Moving Pictures And Audio, MPEG Video Sub-Group, July 2007;\nPaper entitled “A robust content based digital signature for image authentication” by Marc Schneider and Shih-Fu Chang, Columbia University; and\nExtracts from the website of IDDE Inc for Image monitoring service."}
{"text": "Target games which provide a playing piece that is projected toward a target are generally known. In such a game, a player hits a target with a playing piece to score points."}
{"text": "This invention relates to chemical vapor deposition. In particular the invention relates to improved throughput and process cleanliness.\nChemical vapor deposition (CVD) or chemical vapor infiltration (CVI) reactions are processes in which gas phase precursors are reacted together to form a solid phase coating material with resulting gas phase byproducts from the reaction being carried away and trapped or vented. CVD processes may be used in applications that require certain process attributes, including high purity coatings, non-line of sight processing, ability to infiltrate porous substrates and fill open voids with solid phase materials and extremely thin or very tightly controlled coating thickness layers. In most industrial CVD applications, the energy required to cause the reaction between the reactant gasses is supplied in the form of heat in a resistance or induction heated furnace or reactor. Typical CVD reaction temperatures range from 400° C. up to 1600° C. depending on the desired deposition phase, reaction chemistry and substrate material. In applications that require precise coating thickness control or infiltration of a porous substrate, such as the formation of a ceramic matrix composite via CVI, the process is conducted at reduced pressure to reduce the reaction rate and to extend the mean free path of the reactants.\nThe partial pressure of byproducts may increase by several orders of magnitude as they cool and pass through a vacuum system. Depending on the process and byproduct chemistry, the byproducts may stay in the gas phase through the entire vacuum system, or they may condense into a solid or liquid material. In processes where the byproducts tend to condense, great effort is usually used to control where the material condenses so that it can be easily cleaned without dismantling and cleaning the entire CVD system. The condensation of reactant byproducts in vacuum pumps can be particularly problematic causing pump failure in some instances. To circumvent this issue, various cold traps and particle traps are used to condense the byproduct ahead of the pump. Various additional systems may be employed such as heated vacuum lines to insure that byproducts stay in the gas phase until they reach a trap. The presence of reaction byproducts in a vacuum system or process chamber interrupt process yield and adds significant maintenance costs. The problem is particularly acute when the byproduct is hydroscopic and attracts moisture from the atmosphere each time the process chamber is opened for a batch change."}
{"text": "This invention relates to ellipsoidal stage lights and so-called linear moving patterns for projecting patterns from such lights. More particularly, a linear moving pattern attachment is disclosed which enables projection of a pattern having linear motion from such ellipsoidal stage lights.\nEllipsoidal stage lights are old and well known. Specifically, these lights get their name from the ellipsoidal shape of their reflecting mirrors and their use of the two ellipsoidal foci. A brief explanation of this prior art technology is in order.\nAn ellipse drawn on a plane has a major axis and a minor axis. The ellipse is elongate in the direction of the major axis and defines two foci along this major axis. By rotating the ellipse, one obtains an ellipsoid. By cutting off one end of the ellipsoid beyond the second focus, the interior optics of an ellipsoidal stage light can be understood.\nSpecifically, the reflecting mirror of an ellipsoidal stage light is shaped interior of the ellipsoid about one focus of the ellipsoid. The lamp is placed at this focus. When the lamp is energized, light is received and reflected to the remaining second focus of the ellipsoid.\nNaturally, projection to a target, such as a stage must occur. In an ellipsoidal light, the ellipsoid is cut or truncated at the end opposite the reflector, usually slightly beyond the secondary focus. Objective lenses of various kinds and designs are placed in the light beyond the second focus. A beam is then projected by the objective lenses from the light.\nIn a conventional ellipsoidal stage light, shutters are manipulated with respect to the second focus to trim the edges of light projected by the light. Ellipsoidal lights are unique in this ability to use shutters to trim the edges of their projected image; they are simple, inexpensive, and capable of precise edge control of their projected beams in their stage lighting function.\nIn stage or scene light design, it is sometimes necessary to project patterns. It was from this necessity that the so-called xe2x80x9cgoboxe2x80x9d was developed for use with ellipsoidal stage lights.\nVery simply stated, ellipsoidal lights are usually provided with an iris slot at the vicinity of the second focus. By placing a pattern (sometimes called a xe2x80x9cgoboxe2x80x9d), typically in the form of a patterned sheet of metal, glass or temperature resistant film into the iris slot adjacent the second focus of the ellipsoidal light, a corresponding pattern is projected as a light image, usually towards a stage.\nMoving patterns have been projected from such ellipsoidal lights. Unfortunately, the space available at the iris slot in the vicinity of the second focus is extremely limited. To date, only circularly moving patterns have been used with ellipsoidal lights. An example of a gobo having circularly rotating mounted patterns placed in the iris slot of an ellipsoidal light is sold under the trademark GAM TwinSpin II(trademark) of Gam Products Inc. of Hollywood, Calif. In this particular gobo design, circular patterns rotate about the major axis of ellipsoidal light in the vicinity of the second focus. With the GAM TwinSpin II(trademark), it is possible to have two patterns rotate in opposite directions.\nUnfortunately, circularly rotating gobos produce circularly rotating patterns. While such circularly rotating patterns can be occasionally used in light design, they are out of step with the real worldxe2x80x94where most patterns travel either horizontallyxe2x80x94as in the case of cloudsxe2x80x94or verticallyxe2x80x94as in the case of rain (down) or fire (up).\nTo solve this problem for the lighting designer, there has been developed a product known as a projector. An example of such a projector is the GAM Scene Machine sold by Gam Products Inc. of Hollywood, Calif. In short, this projector uses the principle of a light with two foci to project scenes, such as clouds, rain, snow, and fire. Just as in the case of the ellipsoidal light, in this rather expensive dedicated projection instrument, the image to be projected is placed at the second focus of a two focal plane system.\nWhen linearly moving images are desired in such projectors, it is known to use endless pattern loops of continuously rotating metal, or film to project such images. The endless pattern loops typically have a drive roller on one side the projector, a return roller on the opposite side of the projector, and two planar runs of projection material extending in the vicinity of the projectors second focus. Since such projectors are specially designed instruments, space in such instruments is generally of no concern. As a consequence, a relatively large projection rotator is used.\nBy registering one or the other of these planar linearly moving runs of projection material to the second focus of the projector, two separate effects occur. First, the linearly moving run of projection material coincident to the secondary focus has its linearly moving pattern projected by the projector. A corresponding linearly moving pattern of light is projected. Second, the run of projection material not coincident to the secondary focus acts only to attenuate or modulate the light projected. It primarily effects the brightness (or amplitude) of the projected light. Since it is out of the plane of the secondary focus, its image does not project. From such projectors, it is known to project light images of rain, snow, fire, clouds and the like.\nUnfortunately, such projectors are large, specialty instruments. They require special mounting, with their own lighting design and operation techniques.\nIn what follows, we disclose a linear moving pattern for projecting linear motion that is adapted to the narrow confines of a conventional ellipsoidal stage light. In making this adaptation, we set forth the problem to be solved, as well as its solution. The reader will understand that the discovery of the problem to be solved can constitute invention. Therefore, in so far as the prior art does not set forth the problem that we solve, invention is claimed to that aspect.\nA gobo having linear motion pattern projection is inserted to the narrow iris slot of a conventional ellipsoidal stage light. The linear moving pattern includes a housing, a drive roller and motor, a return roller, and an endless pattern loop. The housing has a narrowed end for protruding into the instrument at the iris slot and holding the return roller on one side of the light. The housing has a broad slot end for mounting outside of the ellipsoidal stage light for holding the drive roller and motor at the other side of the light. The broad slot end of the housing has two spring-biased hinged clamps for clamping the gobo into a single unitary linear moving pattern iris insert. A motor clamp holds the drive roller and motor in place. A pattern clamp tensions and dimensions the endless pattern loop for extension into and out of the iris slot of the ellipsoidal light. In assembly, the endless pattern loop is first threaded to the return roller and the return roller mounted to the narrow end of the housing. Second, the drive roller and motor is threaded to the endless pattern loop, placed into the housing at the broad slot end, and locked into place by the motor clamp. Third, the tension clamp is dropped over the endless pattern loop to dimension and tension the endless pattern loop at two parallel and linear runs on opposite sides of the iris slot between the drive roller and return roller. Fourth, the unitary linear moving pattern iris insert is inserted to ellipsoidal stage lighting instrument at the iris slot for operation. In operation, the motor rotates the drive roller relative to the housing to rotate the endless pattern loop. One linear run of the endless loop pattern is registered near the second focus of the instrument to project the desired linear pattern while the remaining linear run of the endless loop pattern attenuates the projected image. There results the ability to provide locally focused and conventional ellipsoidal stage lights with linearly moving projected patterns.\nIt will be understood that one or the other linear runs of the endless pattern loop can be at the focal plane. In this way apparently linear movement of the pattern will be changed (reversed).\nAn advantage of this apparatus and method is that it enables the lighting designer to easily place linearly moving projected patterns wherever an ellipsoidal light can be focused to a light design area. The bulk, added expenses and design inconvenience of specially designed projectors are not required for local linearly moving image projection.\nAn additional advantage of this apparatus is that the linear moving pattern can be easily interchanged among the ellipsoidal lights of a design. The linear moving pattern can be removable between ellipsoidal instruments until an optimal location is determined."}
{"text": "In graphics systems, the graphic display information is contained in a graphics memory at specific locations in the memory. This information is then mapped to the video screen in a one-for-one format. To save time and for convenience, the representations on the screen follow each other sequentially, and the same sequential order is used in memory to store the data for each pixel of information on the screen.\nHowever, problems arise in that video memories have square characteristics with a fixed number of points in the matrix. A video screen, on the other hand, has a number of points called pixels, with each pixel having a number of bits which must be presented to that pixel. Since it is desired to use the same physical memory for many different screen sizes, it is customary to create the memory having a size at least large enough to directly map the largest number of pixels that would be encountered in any one screen. In order to accomplish this goal and not burden the processor with vast numbers of calculations, there must be some easy mathematical coordination between the memory location and the screen location for any data bit. The importance of such ease of calculations can be appreciated when it is realized that in a typical video graphics system each eight bit pixel must be sent to the screen every 12.7 ns. A typical screen would have a pixel array of 1280 by 1024. The display is refreshed 60 times a second. Time spent in processing address information on a per pixel basis then becomes critically important.\nThere are two basic ways to address a pixel in memory. The first of these is the X-Y coordinate method, which seems to be the natural way to think of memory locations. The second method would be to use a linear or vector address giving location data starting from an arbitrary 00 point. Using this system, the processor must calculate the actual position of each pixel.\nProblems arise in the calculations, however, unless special steps are taken to organize the data in the memory exactly as it is presented on the screen. Assume for a moment that the memory is 2,000 columns wide, but the number of necessary screen locations would only take up perhaps 1,500 columns. Visualizing this then, one part of the memory would be vacant. This \"extra\" memory space is hard to use for any other purpose, and thus effectively, wasted.\nIn an attempt to achieve full utilization of the memory, two problems must be solved. One is that there must be a method of removing the information from the ends of the lines of memory and wrapping the end around to the next line of memory. One method of shifting information to a screen is contained in co-pending application entitled \"Graphics Display Split Serial Register System\", Ser. No. 07/387,569, now abandoned, filed concurrently herewith, which application is herein incorporated by reference.\nThe second problem is the restriction that the pixel size must be a power of 2. This restriction stems from the fact that the processor must be able to easily calculate the address of the first pixel in each next row of pixels. The distance between pixel rows is called the pitch. If memory is to be utilized fully, the addresses must be consecutive with the first address in a screen row being the next binary numerical memory address after the address in the preceding row. Since the number of bits per pixel plus the pitch of the screen must be first multiplied and then added together to translate from an X-Y address to a linear address, it follows that any such calculations, because of their great numbers, must be quickly performed. Thus, it is always desired to reduce such calculations to a simple data shift. This can be accomplished when it is realized that the data is binary and thus multiplication by a power of 2 simply requires a one position bit shift, for each such power.\nBased upon the need for quick mathematical operations, a restraint is placed on the number of pixels in a row and this restraint limits the number of pixels to powers of 2.\nAccordingly, a need exists in the art for a system which allows for the utilization of pixels of any size without adding to the processing time for data translation.\nA need also exists in the art for an arrangement which allows the pixel size to be any number and which allows for the packing and shifting of the bits into consecutive memory space so as to conserve memory capacity."}
{"text": "The present invention relates to a processor for processing data of multiple sizes.\nA built-in processor now needs to process data of various sizes, e.g., 8-, 16- and 32-bit data, by itself. However, when applied to a mobile unit, a processor also needs to reduce its power dissipation as much as possible such that the processor can be driven for a longest possible time using a battery with a limited capacity.\nA known RISC (reduced instruction set computer) processor with 32-bit architecture includes: a 32-bit arithmetic and logic unit (ALU) for performing arithmetic and logical operations on data; a 32-bit register file for retaining data therein; and a 32-bit bus for transferring data therethrough. In manipulating 8-bit data, the processor of this type extends the size of the data into 32 bits unconditionally. More specifically, if the 8-bit data in question is signed, then sign extension is carried out on the high-order 24 bits. Alternatively, if the 8-bit data is unsigned, then zero extension is carried out on the high-order 24 bits. Similarly, in manipulating 16-bit data, the processor of this type also extends the size of the data into 32 bits unconditionally. In the conventional processor, the 32-bit data obtained in this manner is provided to the 32-bit ALU and retained in the 32-bit register file or transferred through the 32-bit bus.\nEven in performing a series of operations, like loading two 8-bit data elements from a memory, adding these data elements together and then storing 8-bit data representing the sum in the memory, the RISC processor should use all of its hardware resources, i.e., the 32-bit ALU, 32-bit register file and 32-bit bus. Thus, the processor dissipates power in vain. The same problem arises when the processor handles 16-bit data.\nIt is therefore an object of the present invention to cut down on the power dissipated wastefully by the multi-size data compliant processor.\nTo achieve this object, the inventive processor for handling data of multiple sizes is so constructed as to enable only a part of manipulation means that has a bit width corresponding to a maximum one of those sizes and that is associated with the data size of a specified operand. As used herein, the xe2x80x9cmanipulation meansxe2x80x9d is means for performing an operation, specified by an instruction, on an operand that has also been specified by the same instruction. The manipulation means includes: an ALU for performing arithmetic and logical operations on operand data; a register file for retaining the operand data therein; a bus for transferring the operand data therethrough; a driver for driving the bus; a latch for latching the operand data from the bus; and an extender for extending the size of the operand data. For example, by decoding a given instruction, size information about the data size of the operand specified by the instruction is obtained, thereby enabling only a part of the manipulation means that is specified by the size information.\nTo implement object-oriented architecture, operand data and size information about the size of the operand data are preferably retained in the register file. In this case, if a first instruction decoded requests that operand data of a particular size be written on the register file and if size information about the data size of the operand, specified by the first instruction, has been obtained, then not only the operand data but also the size information are retained in the register file. Alternatively, if a second instruction decoded requests that the operand data retained in the register file be referred to, then size information about the size of the operand data, as well as the operand data itself, are read out from the register file, thereby enabling only a part of the manipulation means that is specified by the size information read out from the register file.\nAccording to another aspect of the present invention, the first instruction may also be an instruction specifying whether the operand data should be handled as signed data or unsigned data. The register file preferably retains the sign information representing whether the operand data is signed or unsigned, in addition to the operand data and the size information about the size of the operand data. In this case, if the sign information representing whether the operand, specified by the first instruction, is signed or unsigned is obtained in accordance with the first instruction decoded, then the register file retains not only the operand data but also the sign information. Alternatively, if the second instruction decoded requests that the operand data retained in the register file be referred to, then the sign information representing whether the operand data is signed or unsigned, as well as the operand data itself, are read out from the register file, and the manipulation means is controlled such that the second instruction is executed in accordance with the sign information read out."}
{"text": "Collapsible canopy frames often include a plurality of telescoping legs, each having one or more X-shaped scissor assemblies extending therebetween. A canopy covering, such as a cloth or leather covering, is disposed above, and supported by, the collapsible canopy frame. The X-shaped scissor assemblies are moveable relative to the telescoping legs to adjust the collapsible canopy frame between an expanded position and a collapsed position.\nIn the expanded position, the collapsible canopy frame provides a temporary shelter. In the collapsed position, the collapsible canopy frame can be more readily transported. Typically, collapsible canopy frames are transported by placing the collapsible canopy frame on a separate wheeled structure, such as a wheeled platform. However, collapsible canopy frames are often relatively heavy and it is therefore desirable to minimize any upward lifting that is required during transportation of a collapsible canopy. Accordingly, a need exists for a wheel assembly and/or components related thereto for attachment to a collapsible canopy frame that facilitates transportation of the collapsible canopy frame."}
{"text": "1. Field Of The Invention\nThe present invention relates generally to liquid crystal display devices, and more particularly to a liquid crystal display device mounted behind a glass cover plate.\n2. Description Of The Prior Art\nIn a liquid crystal cell which serves as an indicator element, the liquid crystal layer is located between two transparent cell plates provided with specific electrode structures. A polarization foil is arranged on both sides of the liquid crystal cell. The rear side of the liquid crystal indicator is provided with a reflector.\nIn appliances with liquid crystal indicator, e.g. measuring instruments, computers, clocks or watches, the liquid crystal display element is conventionally placed in a housing which consists, on the front side, of a cover glass or transparent plastic panel or a glass sheet or sapphire glass. Behind this the liquid crystal display is placed at a distance of some tenths of a millimeter. Even where the liquid crystal display is intended to be in contact with the front glass, a very thin air gap is present between the liquid crystal indicator and the front glass.\nThe use of a front glass in an instrument with a liquid crystal display has the disadvantage, even where the glass possesses good transparency, that reflections and refractions occur at the two surfaces of the front glass and at the surface of the display element, which under unfavorable conditions of illumination, e.g. with weak illumination or with light falling at right angles onto the front glass and dark surroundings, render the reading of the indicated information difficult and reduce the contrast of the indication. With anti-reflection coatings, such as are used for improving optical instruments, it would be possible to eliminate the reflections from the front and rear sides of the cover plate and the surface of the liquid crystal display. This, however, is an expensive and questionable process, especially on the surface of the polarization foil. To make the housing terminate with the liquid crystal display on the front side has the disadvantage that the surface of the liquid crystal display, which is sensitive to scratching, shock and moisture, or the surface of the polarization foil stuck onto the latter, is unprotected."}
{"text": "1. Field of the Invention\nThe present invention relates to a stencil printing apparatus that performs printing using a stencil base sheet (hereinafter referred to also as a “master”) in which perforations are created by platemaking in accordance with the image data of an original mounted on the outer circumferential surface of a plate cylinder, and more particularly to a stencil printing apparatus capable of performing printing on a coated sheet.\n2. Description of the Related Art\nStencil printing apparatuses that comprise a stencil platemaking device (hereinafter referred to simply as a “platemaking” device) for perforating a large number of independent holes in a master using a thermal head to manufacture a master plate, a cylindrical plate cylinder on the outer circumference of which the master made by platemaking is windingly mounted and which is rotationally driven around its axis, an ink supply device provided in the plate cylinder for supplying ink to the inner surface of the plate cylinder, a sheet-feed device for separately feeding out printing sheets, a printing-pressure device for pressing the fed printing sheets onto the outer circumferential surface of the aforementioned plate cylinder, and a discharge device for conveying and discharging the printed sheets, and that employ a system for forming a printed image based on ink having passed through the perforations made by platemaking in the master being transferred onto the printing sheet are already widely known.\nThe large number of independent holes formed by the platemaking device are produced as a result of the perforation of a thermoplastic resin film of the master based on selective heating thereof by very small heat-emitting bodies of a thermal head in accordance with image information. Ink that passes through the perforations formed in the master is directly transferred onto the surface of the printing sheet to form a printed image.\nIn addition, predetermined printing images are formed using this kind of stencil printing apparatus as a result of printing sheets (sheet-like recording medium) separated into individual sheets and conveyed by a sheet-feed portion being pressed, while being conveyed, by the printing-pressure device onto the outer circumferential surface of a printing drum on the outer circumferential surface of which a master manufactured by platemaking is windingly mounted so that ink passing through the perforations of the master is selectively transferred onto the printing sheet surface.\nThe essential conditions required of sheets used in conventional stencil printing apparatuses are a surface of large undulations and ink permeability. This is because the drying system employed in these apparatuses is based on permeation into the sheet. That is to say, rather than natural drying, a system of pseudo-drying the ink transferred onto a sheet based on permeation thereof between the fibers of the sheet and evaporation of the liquid phase of the ink is fabricated.\nAccordingly, the surface of the sheet pressed onto the surface film of the master windingly mounted onto the outer circumference of the printing drum possesses a certain degree of undulations between which an air layer exists.\nBy virtue of this, and the resilient strength of the sheet, a printing sheet can be comparatively easily peeled from the outer circumferential surface of the printing drum immediately following printing by a powerful sheet discharge peeling means. Furthermore, even if the number of printed copies printed using a single master is of the order of 5000 copies and a slight amount of stretch of the master itself has occurred, non-permissible level of creasing and tear will not be generated.\nPrinting using a stencil printing apparatus employing a master used in a platemaking device such as this on a so-called coated sheet on which a coating has been administered on the surface thereof has hitherto hardly even been implemented because it has conventionally been regarded as impossible.\nHowever, there has been a long-held demand in the Japanese domestic market for printing to be able to be performed on a coated sheet. This is because the method employed by manufacturers for the bulk printing of, for example, flyers has concentrated on printing by offset color printing and has involved the flyers being distributed to regional branch shops, and the branch shops printing additional shop map and independent information on a partial blank area of the flyers and distributing them to individual households.\nWhile in this case a lightweight coated sheet is employed as the optimum sheet for offset color printing and, as the optimum printer for the additional printing performed at each branch shop an inexpensive and simple digital platemaking stencil printer able to be used by anyone anytime is employed, regrettably, ink is unable to be used with a coated sheet because it does not dry and fix thereto.\nHowever, stencil printing employing a UV ray-curable ink is already widely known and, if this system is employed, the ink of a printed image can be dried and cured on even non-permeable sheets, and the stencil printing on the aforementioned coated sheet is regarded as being able to be readily performed.\nKnown UV ray-curable ink-employing stencil printing apparatuses are described in Japanese Utility Model Publication No. 4-35188 and Japanese Unexamined Patent Application No. H5-64878.\nIn addition, Japanese Unexamined Patent Application No. 2006-327023 discloses, in a stencil printing apparatus in which both a UV ray-curable ink and a non-UV ray-curable ink, ink attached to a drum being detected as either a UV ray-curable ink or a non-UV ray-curable ink, and a warning being issued when an incompatibility exists. This application similarly discloses a master mounted on a platemaking device being detected as a master suitable for UV ray curing or a master suitable for non-UV ray curing, and a warning being issued when an incompatibility exists.\nIn addition, Japanese Unexamined Patent Application No. 2004-351757 discloses a configuration in which a permeable drying ink and a UV ray-curable ink are selected by switching a supply valve, and which employs a UV ray-curable ink for printing on a coated sheet.\nIn addition, Japanese Unexamined Patent Application Nos. H11-227229 and 2006-168370 disclose, in an ink jet system, detection of a recording medium material type and judgment of whether a first ink or a second ink is to be used in accordance therewith.\nIn addition, Japanese Unexamined Patent Application No. 2001-130120 discloses, in a stencil printing apparatus, a configuration in which two types of master are supported of which one is selectively used in accordance with the printing sheet type.\nIn addition, Japanese Unexamined Patent Application No. 2001-315291 discloses, in a stencil printing apparatus, the adjustment of platemaking energy in accordance with the printing sheet type.\nIn addition, Japanese Unexamined Patent Application No. 2001-328332 discloses, in a stencil printing apparatus, automatic selection of printing conditions in accordance with the printing sheet type.\nIn addition, Japanese Unexamined Patent Application No. 2005-186357 discloses selection of ink jet recording conditions based on a discrimination result of sheet type, and the issue of a warning when an incompatibility exists.\nIn addition, Japanese Unexamined Patent Application No. 2006-281658 discloses, in a stencil printing apparatus, control of UV lamp flashing according to the use of a UV ink drum unit or a non-UV ink drum unit.\nIn addition, Japanese Patent No. 3734247 discloses the discrimination of a coated sheet based on optical detection.\nRepeated implementation of printing tests on coated sheets employing a conventional stencil printing apparatus carried out by the inventor of the present invention with a view to understanding the merits and problems inherent thereto led to an awareness of additional new problems.\nIn stencil printing, a large number of very small perforations are provided by a platemaking process in the film surface of a master windingly mounted on the outer circumference of a printing drum, and ink passing through these perforations is pressed onto the printing sheet by a printing pressure portion and transferred on to the printing sheet surface to form a printed image.\nAt this time, ink is interposed and perfectly closely adhered between the film surface of the master and the surface of the printing sheet. For a coated sheet, because the smoothness of the sheet surface thereof is immeasurably greater than the smoothness of a conventional normal sheet, this adherence force is very large. While the printing sheet is subsequently peeled from the printing drum and discharged, because the printing sheet sticks to the printing drum due to the adhesive strength of the ink, a very large peeling force is required at this time.\nFor a conventional normal sheet, because of its air permeability and the large undulations in the surface thereof, it can be peeled from the printing drum with a moderate level of peeling force. In addition, using a conventional emulsion ink, because both the adhesive strength and the tack value of the ink itself are very low and the ink cohesion force is also low, the ink tears easily and the printing sheet can be peeled off with a moderate level of force.\nHowever, these conditions are significantly different when the printing sheet is a coated sheet. That is to say, because the surface smoothness of a coated sheet is very high and the sheet sticks to the printing drum due to the adhesive strength of the ink, a markedly greater force is required to peel the sheet from the surface of the master film of the outer circumference of the printing drum because. As a result, it was clear that the master film rather than the sheet was lacking in resistance, and that slippage thereof resulting in creasing and tear thereof sometimes occurred.\nIt was also apparent that a switch of the printing ink from a conventional emulsion inks to a UV ray-curable inks of which the principal component is a monomer led to an increase in the adhesion strength and tack value of the ink itself and, in turn, to a magnification of the drawbacks described above.\nAs the master film, a very thin film configured from a thermoplastic synthetic resin of which the thickness thereof allows for the manufacture therein of perforations by platemaking at high speed using a thermal head is used. More specifically, a PET film of thickness of the order of 1.5 to 2 microns is used. A marked drop in strength occurs in the solid image portion in which a large number of perforations are formed by platemaking in this very thin film and, as a result, it separates and creases and tears easily due to the peeling performed immediately following printing."}
{"text": "1. Field of the Invention\nThe invention relates to a process for manufacturing rigid cylindrical cardboard tubes of a certain thickness, whose ends are formed so as to provide ready and efficient closure of the tube, an apparatus for implementing this process and the tubes thus obtained.\n2. Description of the Prior Art\nIt is known from Belgian Pat. No. 435 579 to form paper or flexible cardboard cornets, for example cornets for sweets, whose closure is provided by two facing tongues which are folded back one over the other above the opening of the cornet about a folding line curved towards the point of the cornet. Such cornets are made from a sheet of paper or thin easily deformable cardboard which is cut out when flat and is then shaped so as to form a cornet and is fixed by sticking in this position.\nIt should be understood that the process of Belgian Pat. No. 435 579 cannot be applied to the formation of a cylindrical rigid cardboard tube of a certain thickness, of the kind of tube for example used for sending documents in roll form by post, for it would not be possible starting with rigid cardboard sheets of a certain thickness, cut out when flat, to deform them so as to obtain cylindrical tubes without considerable risk of breaking the cardboard and without creating an unacceptable bulge in the sticking zone.\nThe purpose of the present invention is to overcome these drawbacks and for this it provides a process for forming cylindrical tubes from rigid cardboard, comprising opposite tongues cut out so as to form an hermetic closure system, as well as an apparatus for implementing this process."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the use of genera Prunus endocarp as a dietary supplement product for mammals, to a method for producing the product, to the product produced by the method, and to the method of using the product.\nMore specifically, the dietary supplement is produced from a natural source from which a host of chemical compounds having beneficial properties, as will be described hereinbelow, can be obtained. Also defined hereinbelow is a method for processing this source in a manner which preserves these chemicals and renders the source ingestible.\n2. Description of the Prior Art\nMuch research has been done on the effect of various organic chemical compounds on mammalian bodily function. It is common knowledge that the addition of certain chemical compounds can have beneficial effects. For example, expansive amounts of literature report that phytochemicals of the general class of (a) tannins, (b) phytic acids, (c) phytosterols and (d) flavanoids/isoflavones have (i) astringent properties and (ii) micellar disruption capabilities. As such, they can be used to affect the absorption of fats. Certain of these phytochemicals have also been reported as having ancillary anti-cancer/cancer preventative, hypocholesterolemic, antihepatotoxic and immunostimulant properties as well. Unfortunately, it can be difficult to introduce these compounds into the diet. Synthetic versions of these compounds can be costly as well as unappealing to those who prefer natural products. In addition, even if these compounds can be found naturally in multiple natural sources, there is a practical limit to the number of supplements individuals can reasonably be expected to take. There is therefore a need in the art for a product which provides a multiplicity of such beneficial compounds. Ideally this product would be natural, inexpensive to acquire and process, and easily ingested.\nIt is common knowledge that a large segment of society is overweight. Some have attributed this to a less physically active lifestyle in this modern era; a sedentary lifestyle tends to result in weight gain. Others have blamed consumption of foods with high fat content. Studies have shown that people of excessive weight tend to have a higher incidence of certain health problems. In addition to health risks, society finds excess weight to be aesthetically unappealing.\nAs a result, significant efforts have been directed at finding methods to enable people to lose weight. This has resulted in diet programs under which participants reduce their caloric intake in general and in relation to specific foods which tend to result in increased weight as well as resulting in creation of numerous exercise programs designed to burn calories. In addition, sugar and fat substitutes have been introduced as another way to reduce caloric intake. While some of these efforts may help, they generally prove ineffective for most of the population. Diet and exercise programs are difficult to follow, and many take a dim view with regard to such sugar and fat substitutes.\nA high fat intake diet associated with an inability to match that intake with the expenditure of the consumed fat derived calories is known to be correlated with obesity. It is understood that the body tends to store excess absorbed fat for future use rather than eliminate it, and it is this factor which tends to lead to weight gain. Unfortunately, fats make many food products appealing and this leads to excessive fat consumption. There is therefore a need in the art for a product which reduces the body\"\"s absorption of fats. Ideally this product would be natural and easily ingested. The present invention is designed to satisfy this and other needs.\nAbsorption of fats from the diet in the gut is regulated by gastric secretions from the bile duct. Absorption of the various types of fats typically found in the human diet is controlled by the release of triglyceride digesting esterases which are released from the gallbladder via the bile duct into the upper duodenum. Once in the digestive tract, these esterases hydrolyze the fatty acid chains from the triglycerides, freeing the fatty acids causing the formation of fat micells and consequently allowing fat absorption to take place.\nThe absorption of the free fatty acids from the lumen of the small intestine is achieved through a process of emulsification/micellization with bile salts in the crypts between the microvilli that cover the lumen of the intestine. These created micelles are capable of being absorbed across the cells defining the lumen via a pinocytotic process and into the underlying capillaries that drain into the hepatic portal vein. This diffusion is controlled by the diffusion potential across the brush border of epithelial cells of the micro-villi and the presence of chemical compounds with astringent and/or micellar disruption capabilities that reduce the absorption of fats.\nVolumes of literature report that phytochemicals of the general classes of (a) tannins, (b) phytic acids, (c) phytosterols and (d) flavanoids/isofiavones have (i) astringent properties and (ii) micellar disruption capabilities. As such, they can be used to effectively reduce the absorption of fats. Certain of these phytochemicals are also known to have anti-cancer/cancer preventative, hypocholesterolemic, antihepatotoxic and immunostimulant properties.\nEndocarp preparations from the genera Prunus have been shown to be rich in the phytochemicals referred to above. Absent in the literature, however, is any suggestion of the use of genera Prunus endocarp or portions thereof as a dietary supplement in reducing significantly the absorption of the intake of fats. It is well known that the endocarp of the genera Prunus includes an undesirable compound, hydrocyanic acid. Ideally, this hydrocyanic acid should be removed before genera Prunus endocarp is ingested.\nAccording to the invention there is provided a dietary supplement comprising an endocarp of the genera Prunus in ingestible form.\nFurther according to the invention there is provided a method for using a dietary supplement comprising an endocarp of the genera Prunus including the step of ingesting a sufficient amount of the supplement to reduce the body\"\"s absorption of fat.\nStill further according to the invention there is provided a method for producing a dietary supplement comprising the steps of obtaining a fresh endocarp of the genera Prunus, freezing the obtained structure prior to placing same in a suitable vacuum vessel, connecting the vacuum vessel to a lyophilizer, creating a pressure drop in the vessel to a suitable vacuum level, maintaining the vacuum level for a predetermined duration, returning atmospheric pressure into the vessel, removing the lyophilized structure from within the vessel, and pulverizing the structure into particles of predetermined size.\nYet further according to the invention there is provided a product made by the method of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to an optically induced cell lysis biochip, and more particularly, to a cell lysis method using a system which includes said biochip.\n2. Description of the Related Art\nIn the biomedical field, researches focusing on a specific component within a single cell tend to become more and more sophisticated recently. Therefore, to extract the target proteins, nucleic acids or organelles from the tissues or cells of an organism has also been an important issue. However, traditional tools, including chemical, thermal, electric, ultrasonic, laser, and mechanical approaches, can only lyse a group of cells in a cell-containing solution and further obtain their cell components for study. If one intends to focus the study on the component of a specific cell, he needs to further separate within the group of cells using other methods so as to extract the target components, which is an exhausting and time-consuming process. Moreover, the cell components in the cell-containing solution, such as proteins and nucleus, would easily be denatured if the cell-containing solution is exposed under a non-physiological condition for a long time. Thus, traditional approaches as described above have prevented researchers from quickly extracting specific cell components, causing a research bottleneck in the biomedical field.\nIn view of the above, the present invention has been developed to provide a biochip and a system, with which a single cell can be lysed within a group of cells. Therefore, specific cell components can be further extracted for analysis with less time spent during the process."}
{"text": "It has long been a necessity, particularly within the textile industry, to provide substrates, such as apparel fabrics, as one example, that exhibit a number of simultaneous wash-durable properties. Most notably, water repellency, oil repellency, stain resistance, and stain release characteristics are highly desirable to facilitate cleaning of substrates, if not to prevent complete staining thereof. Unfortunately, provision of such simultaneous and wash-durable characteristics has been severely limited due to the general difficulties with meeting certain surface energy requirements throughout the wash-durable life of such a substrate. Generally, coatings or other treatments have not been readily available or widely known that can provide coexistent water and oil repellency and stain release on a wash durable basis to fabrics (or other surfaces) because the surface energy profile required for one of these properties is disparately different from the surface energy profile required to impart the other property at the same time.\nAlthough there have been some instances of initial simultaneous existence of both properties on certain substrates (as noted below), unfortunately, the degree of wash-durability thereof has been unacceptable for long-term utilization of target substrates. As a result, any significant reduction in either oil or water repellency consequently reduces stain repellency as well. With a reduced propensity to repel stains, the ability to effectuate proper stain release may likewise be diminished, particularly upon exposure to greater degrees of staining and wherein the surface energy profile needed for proper stain release function (which is similar to that needed to impart the aforementioned water and oil repellency properties) is compromised (e.g., is not wash-durable).\nHence, truly effective wash-durable, long-term, stain repellent and stain release treatments have not been forthcoming, since simultaneous prevention of both polar (aqueous) and non-polar (olefinic) liquid penetration into such fabric surfaces has been very difficult to achieve that can withstand extended common laundering procedures. This problem with prior oil and water repellent surface treatments is most prominently observed on typical high stain substrates such as cotton-containing fabrics. Such fabrics are generally difficult to modify at their surfaces to the extent necessary to impart both oil and water repellent features thereto and to retain an acceptable hand. These at least three properties (stain release, water repellency, and oil repellency) are simply unavailable to the textile industry on a wash-durable basis due to the aforementioned surface energy issues. A description of such surface energy properties helps to permit a better understanding of such a phenomenon.\nA fundamental physical property of any material is its surface energy. This property is usually expressed in mJ/m2. Depending on the magnitude of this property, the material may be classified as having a high surface energy or a low surface energy. This property depends generally on the composition of the substrate. For example, a substrate having a surface that contains a significant portion of polar, hydrophilic groups, such as hydroxyl groups, carboxylic acid groups, amine groups, and the like, generally exhibits a high surface energy. Conversely, a substrate having a surface that contains a significant portion of non-polar, hydrophobic groups, such as silicone, fluorinated groups, and the like, generally exhibits a low surface energy. It is readily known that when a polar liquid, such as water, is placed in contact with the surface of a substrate, the liquid will spontaneously wet the surface only if the surface tension of the liquid is lower than the surface energy of the substrate. Conversely, if the surface tension of the liquid is higher than the surface energy of the substrate, spontaneous wetting will not readily occur, and the liquid will remain pooled on the surface of the substrate.\nAs one might expect then, substrate surface energy modification has long been a major field of research for a variety of materials and for a multitude of reasons. For instance, it is often desirable to increase the surface energy of a substrate to facilitate its ability to absorb liquid or to increase the adhesion between a coating and a substrate. Practical examples include the chemical treatment of paper or plastic to enhance their wetting with printing inks and corona treatment of plastic to increase the adhesion between the plastic and another material, such as for the aluminum coating of Mylar® films in packaging applications. Textile substrates have also been modified to create substrates with high surface energy which results in a textile substrate that is hydrophilic and that exhibits improved comfort and stain release properties. As one example, the detergent industry has employed this technique for determining effective methods of cleaning various textile substrates.\nSurface energy modification has also been utilized in other coating applications, such as to produce non-stick surfaces exhibiting low surface energy through the application of Teflon™ to cookware and cooking utensils. Textile substrates have also been modified with low surface energy treatments in order to produce textile substrates that are hydrophobic and that exhibit repellent properties (such as for water repellent rainwear).\nIt has commonly been observed that substrates treated with fluorinated polymers generally exhibit a contact angle of greater than 100 degrees with water. The advancing and receding contact angles are very similar. The major component of the surface energy of such treatments is dispersive. Substrates treated with dual functional repellents, such as disclosed in U.S. Pat. No. 3,574,791 to Sherman et al., generally exhibit lower contact angles with water when compared with traditional fluorochemical repellents, and therefore, tend to exhibit lower repellency. The measured surface energy contains significant dispersive and polar components. Differences can usually be measured between the advancing and receding contact angles.\nIn some instances, a measurable degree of hysteresis exists between the advancing and receding contact angle, indicating that the surface energy has changed in the presence of a liquid. Barring liquid adsorption, hysteresis is indicative that the surface energy has changed (kinetically or thermodynamically) in the presence of a liquid or environmental condition. This measurable degree of hysteresis provides further evidence that the substrate is autoadapting to its environment. One method for achieving ideal performance for textile applications would be obtained from a composition that provides high advancing contact angles (i.e., >90 degrees), exhibiting non-porous behavior, to impart stain resistance and provides low receding contact angles (i.e., <90 degrees), exhibiting porous behavior, to impart stain release to the substrate. Another method to achieve ideal performance for such applications would be obtained from a composition that imparts high advancing and high receding contact angles between a staining substance and the substrate, followed by low advancing and receding contact angles during exposure to a cleaning procedure.\nIt would be desirable for a porous or stainable surface to exhibit high contact angles versus a variety of liquids to prevent adsorption or staining. It would also be desirable for such surfaces to adapt to a change in their environment, such as in a cleaning medium, to enhance removal of stains and soil. Other environmental conditions that could induce a change in the surface energy of a substrate include changes in temperature, moisture content, and other environmental factors. Highly desirable would be a surface that reversibly adapts to its environment, such that the surface is stain resistant and cleanable and retains this effect through a number of use cycles. In many end-use applications such as apparel, carpet, upholstery, and the like, appearance retention of the product is extremely important. While stain resistant treatments have been developed for each of these exemplary applications, it has been found, that much like stain resistant apparel treatments, such treatments have an adverse effect on subsequent cleaning. Thus, it would be highly desirable to develop soil and stain resistant textile substrates, regardless of the end-use application, that possess enhanced cleanability using appropriate cleaning techniques.\nWith the development of XPS, SIMS, and other surface analytical techniques, it has become possible to detect certain chemical groups at the surface of materials. For instance, one can measure the concentration and depth profile of functional groups, such as CF3 moieties commonly found in fluoropolymer stain resist chemicals. Through appropriate sample preparation techniques, it is also possible to observe changes that take place on the surface of a substrate and that occur as a result of changes in the environment to which the substrate is exposed. For example, a substrate that is observed to contain predominately low surface energy groups, such as CF3 groups, under a first set of conditions can be shown to contain significant hydrophilic high surface energy groups, such as hydroxyl groups, at its surface under a different, second set of conditions. This polarity change typically allows the surface of the substrate to wet (i.e., absorb liquid), thereby enhancing stain release. As the substrate's environment is returned to the first set of conditions, one can observe, for example, the CF3 groups return to the substrate's surface, thus, returning the substrate to its low surface energy, stain resistant state.\nSome treatment compositions, such as polymers, possess other properties, such as glass transition temperature, which may influence the ultimate performance of the treated substrate. For instance a hard polymer that is characterized by a high glass transition temperature may provide increased protection against wetting, especially forcibly wetting. However, this stiff, high glass transition polymer would likely require more work to adapt to changes in its environment due to less intra-polymer flexibility. In addition, the polymer molecular weight and addition of co-monomers may enhance wetting, adhesion, chemical reactivity, and durability for a variety of substrates as well.\nAs should thus be evident, modification to provide a proper surface energy profile to impart simultaneous wash-durable oil repellency, water repellency, stain resistance, and stain release properties to a target substrate has been sought after for many years without success.\nThe invention as described herein illustrates that certain combinations of chemicals and processing conditions permit and/or facilitate tailoring of the surface properties of a target substrate to obtain the desired balance of surface energy profiles to impart simultaneous repellency and stain release characteristics thereto. Furthermore, this unique combination of features has surprisingly been shown to be quite durable upon exposure to routine as well as industrial cleaning methods."}
{"text": "Malignant tumors often express characteristic antigens or \"markers\" which offer a mechanism for tumor prevention, resistance or treatment. The antigens which are characteristic of the tumor may be purified and formulated into vaccines. This may stimulate an antibody response and a cellular immune response which are helpful in controlling tumor growth. At a minimum, the antibodies raised by these antigens can be used as detection tools to monitor the level of tumor marker in the host to track the course of the disease or the effectiveness of treatment.\nIt is well known that the immunogenicity of antigens can be enhanced by coupling these hapten-bearing moieties to carriers. A variety of carriers are routinely used, such as keyhole limpet hemocyanin and various serum albumins. It is also understood that certain cytokines, such as GM-CSF, have the capacity to enhance primary antibody responses to antigens (Morrissey, P. J., et al, J Immunol (1987) 139: 1113-1119).\nThe immune response-enhancing ability of a cytokine, when coupled to a viral antigen, has been shown by Hinuma, S., et al, FEBS (1991) 288: 138-142. In this work, interleukin-2 was coupled to a herpes simplex virus type I glycoprotein by generating a fusion protein consisting of the glycoprotein D and human IL-2. The conjugate was shown to induce high antibody responses and cell-mediated immunity to HSV-I in mice.\nIt has now been found that B cell lymphoma tumor-associated antigens can be coupled to immune response-enhancing cytokines, such as GM-CSF, IL-2 and IL-4, to produce an immune response to the tumor antigen and enhance the ability of the host to resist tumor growth associated with the antigen."}
{"text": "This invention relates to improvements to suspending loads from a rotary wing aircraft (preferably a helicopter) for applications such as placing a human upon an electric transmission line, or transmission line tower, wherein the transmission line may be either energised (live) or de-energised (dead).\nConstructing, maintaining and repairing electric transmission line systems are done either in an energised or de-energised state. When repairing or maintaining such systems it is generally uneconomical to de-energise the system. Hence, it is preferred that maintenance or repair be carried out upon energised conductors, Where this is not possible maintenance and repairs are required to be done as quickly as possible so that loss of revenue and possible inconvenience to users is minimised.\nOne approach to repair and maintenance has been that of humans climbing the large transmission system towers to access conductors, insulators, and the various pars of the towers themselves. This is a time consuming and inefficient approach in which vast amounts of revenue may be lost in cases when the lines have to be de-energised for safety reasons.\nAnother approach is that of using both insulated and non-insulated vehicles or towers in which when working upon energised lines such equipment must be insulated. Nevertheless, there is still a potential for fault discharges due to structural defaults or negligence of personnel which can cause damage to equipment and more seriously death to personnel. Another disadvantage is that this equipment is not suited to rugged terrain and its transportation speed is relatively slow. Consequently, when considering Australian transmission line distances and terrain, repair and maintenance may be time consuming and impractical using this approach.\nAnother approach is that of carrying personnel and equipment upon a helicopter wherein a human sits on a rigid platform attached to the helicopter. The helicopter manoeuvres close to a desired area (the area being either energised or de-energised) and the worker then carries out maintenance, repair, or installation whilst being supported by the helicopter. It is possible for humans to work on energised conductors using this approach as both the human and helicopter are ungrounded and therefore in an ideal situation there is no danger of damage or death. Unfortunately, this is a relatively expensive approach as there is required at least a helicopter pilot, a service and maintenance engineer, and a helicopter which must be in operation throughout the maintenance or service operation. In addition, there is a danger that the helicopter's rotor blades may cause a short between two conductors or may hit part of the transmission system which may just damage the system or could snap a rotor blade causing death to the pilot and engineer.\nAnother more suitable approach has been developed and published in U.S. Pat. No. 4,673,059 in which a maintenance engineer is suspended from an ungrounded support means such as a helicopter. The engineer may be suspended from a suspension means of insulated components so that the probability of shorting between conductors or a grounding discharge is further reduced. The suspension means is connected to a cargo hook attached to the helicopter and the engineer sits in a basket or seat device. There is a limited slip means to allow for a margin of error when the load (engineer) is simultaneously tethered to the helicopter and stationary component of the transmission line system. As an alternative to the limited slip device there is a common break-away link adapted to rupture when an instantaneous dynamic force of about 900 lbs is applied.\nUnfortunately, the invention disclosed in U.S. Pat. No. 4,673,059 has problems associated with safety of suspended loads (especially when near live transmission lines). One such problem with the above arrangement is that in one embodiment the load is suspended from only a cargo hook in which there is no backup safety feature. Thus, the pilot may inadvertently or erroneously release the hook therefore endangering the engineer and other people."}
{"text": "In the gaming industry, certain games require that batches of randomly shuffled cards are provided to players and sometimes to dealers in live card games. It is important that the cards are shuffled thoroughly and randomly to prevent players from having an advantage by knowing the position of specific cards or groups of cards in the final arrangement of cards delivered in the play of the game. At the same time, it is advantageous to have the deck(s) shuffled in a very short period of time so that there is minimal downtime in the play of the game.\nBreeding et al., U.S. Pat. Nos. 6,139,014 and 6,068,258 (assigned to Shuffle Master, Inc.) describe a machine for shuffling multiple decks of playing cards in a batch-type process. The device includes a first vertically extending magazine for holding a stack of unshuffled playing cards, and second and third vertically extending magazines each for holding a stack of cards, the second and third magazines being horizontally spaced from and adjacent to the first magazine. A first card mover is positioned at the top of the first magazine for moving cards from the top of the stack of cards in the first magazine to the second and third magazines to cut the stack of unshuffled playing cards into two unshuffled stacks. Second and third card movers are at the top of the second and third magazines, respectively, for randomly moving cards from the top of the stack of cards in the second and third magazines, respectively, back to the first magazine, thereby interleaving the cards to form a vertically registered stack of shuffled cards in the first magazine. Elevators are provided in the magazines to bring the cards into contact with the card movers. This shuffler design is currently marketed under the name MD1® shuffler and MD1.1™ shuffler in the United States and abroad.\nSines et al., U.S. Pat. No. 6,019,368 describes a playing card shuffler having an unshuffled stack holder that holds an in-feed array of playing cards. One or more ejectors are mounted adjacent the unshuffled stack holder to eject cards from the in-feed array at various random positions. Multiple ejectors are preferably mounted on a movable carriage. Extractors are advantageously used to assist in removing playing cards from the in-feed array. Removal resistors are used to provide counteracting forces resisting displacement of cards, to thereby provide more selective ejection of cards from the in-feed array. The automated playing card shuffler comprises a frame; an unshuffled stack holder for holding an unshuffled array of playing cards in a stacked configuration with adjacent cards in physical contact with each other and forming an unshuffled stack; a shuffled array receiver for holding a shuffled array of playing cards; at least one ejector for ejecting playing cards located at different positions within the unshuffled stack; and a drive which is controllable to achieve a plurality of different relative positions between the unshuffled stack holder and the at least one ejector. This shuffler design is currently marketed under the name RANDOM EJECTION SHUFFLER™.\nGrauzer et al., U.S. Pat. No. 6,149,154 (assigned to Shuffle Master, Inc.) describes an apparatus for moving playing cards from a first group of cards into plural groups, each of said plural groups containing a random arrangement of cards, said apparatus comprising: a card receiver for receiving the first group of unshuffled cards; a single stack of card receiving compartments generally adjacent to the card receiver, said stack generally adjacent to and movable with respect to the first group of cards; and a drive mechanism that moves the stack by means of translation relative to the first group of unshuffled cards; a card moving mechanism between the card receiver and the stack; and a processing unit that controls the card moving mechanism and the drive mechanism so that a selected quantity of cards is moved into a selected number of compartments. This shuffler is currently marketed under the name ACE® shuffler in the United States and abroad.\nGrauzer et al., U.S. Pat. No. 6,254,096 (assigned to Shuffle Master, Inc.) describes an apparatus for continuously shuffling playing cards, said apparatus comprising: a card receiver for receiving a first group of cards; a single stack of card receiving compartments generally adjacent to the card receiver, said stack generally vertically movable, wherein the compartments translate substantially vertically, and means for moving the stack; a card moving mechanism located between the card receiver and the stack; a processing unit that controls the card moving mechanism and the means for moving the stack so that cards placed in the card receiver are moved into selected compartments; a second card receiver for receiving cards from the compartments; and a second card moving mechanism between the compartments and the second card receiver for moving cards from the compartments to the second card receiver. This shuffler design is marketed under the name KING® shuffler in the United States and abroad.\nJohnson et al., U.S. Pat. No. 5,944,310 describes a card handling apparatus comprising: a loading station for receiving cards to be shuffled; a chamber to receive a main stack of cards; delivery means for delivering individual cards from the loading station to the chamber; a dispensing station to dispense individual cards for a card game; transfer means for transferring a lowermost card from the main stack to the dispensing station; and a dispensing sensor for sensing one of the presence and absence of a card in the dispensing station. The dispensing sensor is coupled to the transfer means to cause a transfer of a card to the dispensing station when an absence of a card in the dispensing station is sensed by the dispensing sensor. Individual cards delivered from the loading station are randomly inserted by an insertion means into different randomly selected positions in the main stack to obtain a randomly shuffled main stack from which cards are individually dispensed. The insertion means includes vertically adjustable gripping means to separate the main stack into two spaced apart substacks to enable insertion of a card between the substacks by the insertion means. The gripping means is vertically positionable along the edges of the main stack. After gripping, the top portion of the stack is lifted, forming two substacks. At this time, a gap is created between the stacks. This shuffler is marketed under the name QUICKDRAW™ shuffler in the United States and abroad.\nSimilarly, Johnson et al., U.S. Pat. No. 5,683,085 describes an apparatus for shuffling or handling a batch of cards including a chamber in which a main stack of cards is supported, a loading station for holding a secondary stack of cards, and a card separating mechanism for separating cards at a series of positions along the main stack. The separating mechanism allows the introduction of cards from the secondary stack into the main stack at those positions. The separating mechanism grips cards at the series of positions along the stack and lifts those cards at and above the separation mechanism to define spaces in the main stack for introduction of cards from the secondary stack. This technology is also incorporated into the QUICKDRAW™ product.\nSines et al., U.S. Pat. No. 5,676,372 describes an automated playing card shuffler, comprising: a frame; an unshuffled stack holder for holding an unshuffled stack of playing cards; a shuffled stack receiver for holding a shuffled stack of playing cards; at least one ejector carriage mounted adjacent to said unshuffled stack holder, said at least one ejector carriage and said unshuffled stack holder mounted to provide relative movement between said unshuffled stack holder and said at least one ejector carriage; a plurality of ejectors mounted upon said at least one ejector carriage adjacent the unshuffled stack holder for ejecting playing cards from the unshuffled stack, the ejecting occurring at various random positions along the unshuffled stack.\nJohnson et al., U.S. Pat. No. 6,267,248 describes an apparatus for arranging playing cards in a desired order, said apparatus including: a housing; a sensor to sense playing cards prior to arranging; a feeder for feeding said playing cards sequentially past the sensor; a storage assembly having a plurality of storage locations in which playing cards may be arranged in groups in a desired order, wherein the storage assembly is adapted for movement in at least two directions during shuffling; a selectively programmable computer coupled to said sensor and to said storage assembly to assemble in said storage assembly groups of playing cards in a desired order; a delivery mechanism for selectively delivering playing cards located in selected storage locations of the storage assembly; and a collector for collecting arranged groups of playing cards. The storage assembly in one example of the invention is a carousel containing a plurality of card storage compartments. The device describes card value reading capability and irregular (e.g., missing or extra) card indication. The desired orders described include pack order and random order.\nGrauzer et al., U.S. Pat. No. 6,651,981, assigned to Shuffle Master, Inc., describes a device for forming a random set of playing cards including a top surface and a bottom surface, and a card receiving area for receiving an initial set of playing cards. A randomizing system is provided for randomizing the initial set of playing cards. A card collection surface is located in a card collection area for receiving randomized playing cards, the card collection surface receiving cards so that all cards are received below the top surface of the device. An elevator is provided for raising the card collection surface so that at least some randomized cards are elevated at least to the top surface of the device. A system for picking up segments of stacks and inserting cards into a gap created by lifting the stack is described.\nU.S. Pat. No. 5,605,334 to McCrea, Jr., describes a secure game table system for monitoring each hand in a progressive live card game. The progressive live card game has at least one deck with a predetermined number of cards, the secure game table system having players at a plurality of player positions and a dealer at a dealer position. The secure game table system comprises: a shoe for holding each card from at least one deck before being dealt by the dealer in the hand, the shoe having a detector for reading at least the value and the suit of each card, the detector issuing a signal corresponding at least to the value and suit for the card. A card-mixing system may be combined or associated with the card-reading shoe. A progressive bet sensor is located near each of the plurality of player positions for sensing the presence of a progressive bet. When the progressive bet is sensed, the progressive bet sensor issues a signal corresponding to the presence of the wager. A card sensor located near each player position and the dealer position issues a signal when a card in the hand is received at the card sensor. A game control has a memory and is receptive of progressive bet signals from the progressive bet sensor at each player position for storing in memory which player positions placed a progressive bet. The game control is receptive of value and suit signals from the detector in the shoe for storing in memory at least the value and suit of each card dealt from the shoe in the hand. The game control is receptive of card-received signals from card sensors at each player position and the dealer position for correlating in memory each card dealt from the shoe in game sequence to each card received at a player position having a progressive bet sensed. The specification indicates that FIG. 16 is an illustration setting forth the addition of a single card reader to the automatic shuffler of U.S. Pat. No. 5,356,145, to Verschoor. In FIGS. 16 and 17 is set forth another embodiment of the secure shuffler of the U.S. Pat. No. 5,605,334, based upon the shuffler illustrated in FIGS. 12-16 of U.S. Pat. No. 5,356,145. The shuffler may be mounted on a base in which is contained a camera with a lens or lenses and the camera may be embedded in a base of the shuffler.\nU.S. Pat. No. 6,361,044 to Block et al. describes a top of a card table with a card-dispensing hole therethrough and an arcuate edge covered by a transparent dome-shaped cover. A dealer position is centrally located on the tabletop. Multiple player stations are evenly spaced along the arcuate edge. A rotatable card placement assembly includes an extendable arm that is connected to a card carrier that is operable to carry a card. In response to signals from a computer, the rotation of the assembly and the extension of the arm cause the card carrier to carry the card from the card dispensing hole to either the dealer position or any of the player positions. The card carries a barcode identification thereon. A barcode reader of the card carrier provides a signal representation of the identification of the card to the computer.\nU.S. Pat. No. 6,403,908 to Stardust et al. describes an automated method and apparatus for sequencing and/or inspecting decks of playing cards. The method and apparatus utilize pattern recognition technology or other image comparison technology to compare one or more images of a card with memory containing known images of a complete deck of playing cards to identify each card as it passes through the apparatus. Once the card is identified, it is temporarily stored in a location corresponding to or identified according to its position in a properly sequenced deck of playing cards. Once a full set of cards has been stored, the cards are released in proper sequence to a completed deck hopper. The method and apparatus also include an operator interface capable of displaying a magnified version of potential defects or problem areas contained on a card which may be then viewed by the operator on a monitor or screen and either accepted or rejected via operator input. The device is also capable of providing an overall wear rating for each deck of playing cards.\nMany other patents provide for card reading capability in different physical manners, at different locations, and in different types of apparatus, from card reading shoes, to card reading racks, to table security control systems, such as are disclosed in U.S. Pat. Nos. 4,667,959 (Pfeiffer et al.), 6,460,848 (Soltys et al., assigned to MindPlay LLC), 6,270,404 (Sines et al., automated system), 6,217,447 (Lofink et al.), 6,165,069 (Sines et al.), 5,779,546 (Meissner et al.), 6,117,012 (McCrea, Jr.), 6,361,044 (Block et al.), 6,250,632 (Albrecht), 6,403,908 (Stardust et al.), 5,681,039 (Miller), 5,669,816 (Garczynski et al., assigned to Peripheral Dynamics), 5,722,893 (Hill et al., assigned to Smart Shoes, Inc.), 5,772,505 (Garczynski et al., assigned to Peripheral Dynamics), 6,039,650 (Hill, assigned to Smart Shoes, Inc.), 6,126,166 (Larson et al., assigned to Advanced Casino Technologies), and 5,941,769 (Order, unassigned).\nU.S. Pat. No. 6,629,894 to Purton and assigned to Dolphin Advanced Technologies Pty Ltd., of Victoria, Australia, discloses an apparatus for verifying a deck or plural decks of cards. The device includes a card in-feed tray, a card moving mechanism, a camera, a processor located on a card transport path and an accumulation tray. The apparatus is incapable of shuffling cards. Cards can be fed from either tray past a camera in order to verify the deck. The processor compares the read cards with stored values and then reports that outline deviations from expected values are printed. Examples of printed reports include the rank and suit of each card that is missing, or the rank and suit of extra cards present.\nAlthough these and other structures are available for the manufacture of playing card shuffling apparatus, new improvements and new designs are desirable. In particular, it would be desirable to provide a batch-style shuffler that is faster, provides random shuffling, is more compact than currently available shuffler designs and is capable of reading the rank and/or suit of each card. Additionally, it would be desirable to use the device of the present invention to verify decks of cards either prior to use or as part of the decommissioning process."}
{"text": "Given the ever more stringent environmental requirements being placed on motor vehicles, efforts are underway to reduce fuel consumption, and hence CO2 emissions, to the greatest extent possible. One approach is to give the motor vehicle as lightweight a design as possible. To this end, increasing use is being made of light metal components. These components are frequently fastened to the vehicle body via adhesive bonding, with the latter normally being made out of steel.\nDE 103 600 350 A1 describes a problem that arises therefrom. It is the result of the various painting and subsequent drying processes that a motor vehicle must go through during its manufacture. In motor vehicles with a steel body and a roof made out of a light metal, the adhesive bond may experience warpage, for example bulging effects in the vehicle roof, which are caused by the differing thermal expansion behavior of the materials.\nDE 103 60 350 A1 proposes a method for manufacturing a motor vehicle in which the production processes leading up to a finished adhesive bond between a roof module and a body module can be controlled in such a way as to preclude warpage owing to differing thermal expansion behavior to the greatest possible extent."}
{"text": "The present invention relates to linseed oil which has been processed for the purpose of improving the properties of the oil in connection with the use in for example impregnation of products based on cellulose, such as wood. The invention also relates to a process for the manufacture of such linseed oil.\nEven if the present invention is applicable in relation to other cellulose-based products than wood the invention will in the following be described in connection with wood in the form of different types of lumber. Decomposition of wood when used outdoors is almost exclusively caused by fungi and bacteria. The requirement for growth of such micro-organisms is the presence of water. Furthermore, a suitable temperature and access to nutrition and oxygen are required. Prevention of decomposition can be based on the elimination of one or more of these conditions necessary for the growth of the micro-organisms. As an example there can be mentioned that if the moisture content is brought down to a value lower than about 20% biological decomposition is substantially prevented.\nA number of different impregnation processes for the protection of wood are known, and these methods can in principle be divided up into techniques based on the supply of a substance poisonous to the growth of the micro-organisms, and techniques whereby the wood is given hydrofobic properties, for example by impregnation with some kind or oil. In the latter case an impregnation agent is crude linseed which can be of hot- or cold-pressed type, both of which by suitable heat treatment are concerted into boiled linseed oils. The use of such linseed oils is, however, associated with certain drawbacks for example a viscosity which is unfavourably high for ease of handling, a contents of components constituting substrate for the relevant micro-organisms, unfavourable drying properties, etc. One method used for lowering the viscosity is dilution with a solvent something which, however, is environmentally unacceptable. The use of chemical poisons is presently frequently used but should for the future be mainly prohibited for environmental reasons. Also impregnation with a combination of hydrofobation and toxic treatment is used, for example treatment with creosote.\nThe present invention has for an object to provide a processed linseed oil with properties which are substantially improved, mainly in connection with the use of the linseed oil for the impregnation of wood and wood products.\nAnother object of the invention is to provide a process for the manufacture of such linseed oil of improved properties.\nYet an object of the invention is to provide a process, whereby the uptake of the oil by the wood or the wood product in connection with impregnation is substantially increased and can even reach a near 100 percent level of uptake.\nIn connection with extensive research and experimental work it has in accordance with the invention been found that a treatment of linseed oil resulting a situation that its content of free tocoferol will be less than about 100 ppm, results in substantially improved properties of the linseed oil, particularly in connection with its use as an impregnating agent for wood or lumber.\nThese and other objects of the invention which will be clear from the following description are obtained by a processed linseed oil essentially consisting of linolenic, linolic and oleic acids, mainly in the form of triglycerides, the linseed oil being characterized in that its contents of free tocoferol is less than about 100 ppm.\nIt is particularly preferred that the processed linseed oil contains free tocoferol in a concentration less than about 75 ppm and particularly less than about 50 ppm.\nEven if the processed linseed oil according to the present invention can be based on cold pressed as well as hot pressed linseed oil it is preferred to use cold-pressed linseed oil as a starting material for the linseed oil according to the invention.\nThe new process according to the invention for the manufacture of a linseed oil of the above type is characterized by the following steps:\na) heating crude linseed oil to an increased temperature lower than the boiling point of water,\nb) adding a heated inorganic acid to the oil and mixing the oil,\nc) separation of precipitated materials, for example by sedimentation and discharge and removal of bottom deposit formed,\nd) adding a heated aqueous solution of alkali to the oil and mixing the oil,\ne) discharge of the soap formed, and\nf) rinsing the oil with hot water until a clear rinsing water is obtained.\nThe heating in step a) above suitably takes place at a temperature lying within the range about 80 to about 99xc2x0 C. while avoiding that the boiling point of water will be reached.\nThe inorganic acid is preferably added in an amount of about 0.5 to 30% by weight based on the weight of the oil. Among suitable inorganic acids there may be mentioned phosphoric acid, sulphuric acid and hydrochloric acid. It is preferred to use as an inorganic acid phosphoric acid, for example thermal phosphoric acid, at a concentration of for example 80% and in a quantity of about 0.5 to about 2% by weight, particulary from about 1 to about 1.5% by weight.\nThe mixing in step b) suitably lasts until the reaction has come to a stop, which corresponds to a period of about 10 to about 30 minutes.\nAs alkali any suitable alkaline compound can be used, for example hydroxides of alkali or earth alkali metals, particularly sodium hydroxide or potassium hydroxide. Particularly preferred is sodium hydroxide, also called caustic soda. The alkaline compound is suitably added in a quantity equivalent to about 2 to 5% by dry weight based on the weight of the oil. Alkali is suitably added in the form of an aqueous solution in a quantity of about 7 to 15% by weight water also this based on the weight of the oil.\nFor the purpose of avoiding thickening of the oil during heating the oil may in a further step be heated while adding a suitable agent, such as manganese tetraborate. A suitable temperature range for the heating of the oil in this further step is from about 120 to about 140xc2x0 C. In order to improve the trying properties of the processed oil a suitable catalyst can be added, for example example manganese tetraborate. The agent used as a catalyst can be the same as the agent used to avoid thickening of the oil during heating.\nThe invention will in the following be further described with reference to non-limiting examples, wherein quantities and percentages relate to weight if not otherwise stated. In the examples abbreviations concerning the fatty acid present in the linseed oil are furthermore used. Accordingly, for example oleic acid is abbreviated C18:1, linolic acid C18:2 and linolenic acid C18:3. The number following the colon thus indicates the degree of unsaturation."}
{"text": "1. Field of the Invention\nThe present invention relates to an improved apparatus and process for extracting high-purity aromatics from gasoline which increases system capacity while reducing energy consumption.\n2. Prior Art\nSeveral commercially proven processes and apparatus are available for extracting high-purity aromatics from gasoline, coke oven light oil, and pyrolysis naphtha. Most include a liquid-liquid extractor followed by an extractive distillation column for extracting a high-purity aromatic stream, apparatus for removing solvent from the product streams, and solvent conditioning facilities.\nThe extraction of benzene and heavier aromatic homologs has been practiced commercially for about a century. Prior to the preparation of high-octane gasoline from crude oil, aromatics were extracted from liquids produced during the coking and gasification of coal. With the advent of platinum reforming (xe2x80x9cPlatformingxe2x80x9d) in the late 1940\"\"s, a large source of less expensive aromatics became available in oil refineries.\nAt about the same time Dow and others were developing commercial plants to produce ethylene glycol for the automotive antifreeze market. One of the heavy byproducts of this process was di-ethylene glycol. Dow found that this material could be used to extract aromatics from gasoline, and developed a process to accomplish this.\nDow made an arrangement for UOP to market the process once it was proved, naming it UDEX in honor of the new partnership promoting the process. This process dominated the extraction field through the 1950\"\"s until the Shell Sulfolane process supplanted it in the 1960\"\"s.\nThe early UDEX units used di-ethylene glycol (xe2x80x9cDEGxe2x80x9d) and di-glycol amine (xe2x80x9cDGAxe2x80x9d) for solvents. The consumption of energy was typically in the range of 1200 to 1500 BTU/pound of extract. In the early 1960\"\"s, tri-ethylene glycol replaced most of the DEG/DGA, reducing energy consumption to 1000 to 1200 BTU/pound of extract, and increasing unit capacity by 20 to 30%. In the 1970\"\"s, tetra-ethylene glycol replaced most of the tri-ethylene. With this change, the energy consumption dropped to the range of 800 to 1000 BTU/pound of extract and the capacity increased another 10 to 20%. A solvent additive called xe2x80x9cCaromxe2x80x9d was used in the 1960\"\"s to decrease the energy consumption and increase capacity, each changing in the range of 5 to 10%.\nThe introduction of the Shell Sulfolane process in the 1960\"\"s ended the design and construction of most UDEX apparatus. The Shell Sulfolane apparatus usually consumes less than 700 BTU/pound of extract. While this is a strong advantage, the process has two important disadvantages.\nFirst, the Sulfolane process requires four large columns rather than the two required for the UDEX process. This increases capital cost.\nSecond, the solvent can become corrosive. Reboiler replacement and exotic metallurgy are not uncommon due to this corrosive nature. Entire columns have had to be replaced at times.\nThus, because of the low capital cost and non-corrosive nature of the glycol units, a UDEX apparatus having a low consumption of energy would have substantial application in the aromatics field.\nAccordingly, it is a primary object of the invention to produce an improved process and apparatus for glycol based extraction of aromatics from gasoline and the like commonly referred to as the UDEX process.\nSuch object, as well as others, is accomplished by the process and apparatus of the present invention which decrease liquid-vapor flashing, reduce reflux flow rate, and use heat of enthalpy produced at one point as a source of energy used at another point, decreasing energy consumption while significantly increasing purity and amount of product obtained."}
{"text": "1. Field Of The Invention\nThis invention relates to the art of the electroluminescence and more particularly, to an inorganic thin film electroluminescence device (hereinafter referred to simply as EL device) using inorganic fluorescent materials, which device is particularly suitable for used as a flat light source or a flat panel display.\n2. Description Of The Prior Art\nInorganic thin film EL devices have recently drawn attention to their use as a flat panel display or flat light source. For this purpose, there has been heretofore used inorganic fluorescent materials, which comprise a matrix made of at least one of ZnS, CaS, SrS and the like and doped with not larger than 5 wt % of at least one element selected from Mn, Tb, Sn, Ce, Eu, Sm, Tm and the like as a light emission center.\nAmong these fluorescent materials, ZnS:Mn has been studied for use as an orange color fluorescent material. Because of its good luminance, life and the like characteristics, this fluorescent material has been utilized as a flat panel display. On the other hand, lead sulfides such as ZnS:Tb have been widely studied for use as a green color fluorescent material.\nOther matrices which have been extensively studied include alkaline earth metal sulfide fluorescent material such as SrS:Ce as a blue light emitting fluorescent material, CaS:Eu as a red light emitting fluorescent material, a green light emitting fluorescent material, and the like.\nThe light emitting mechanisms of the fluorescent material or fluorescer of transition metals and rare earth elements are different from each other: with transition metals such as Mn, the light emission is caused by direct collision of the electrons of the matrix, whereas with rare earth elements, light emission depends greatly on the transition of an energy corresponding to the band gap of the matrix. The alkaline earth metal sulfides have a band gap energy ranging from 4.3 to 4.4 eV and ZnS has a band gap energy as small as 3.6 eV. In order to obtain blue to UV light emission which requires a high energy, these band gap energies are too small. Accordingly, for the blue to UV light emission, studies have been made on matrices which have a high band gap energy and include, for example, ZnF.sub.2 :Gd whose band gap energy ranges 7 to 8 eV (J. J. A. P. vol. 10 B (1991), pp. L1815 to 11816) and CaF.sub.2 :Eu (Appl. Phys. Lett. 41, 1982, p.462).\nInorganic fluorescers using these sulfides except ZnS:Mn are not always satisfactory for use as a flat panel display or flat light source with respect to the luminance, efficiency and life, and have never been applied as a flat panel color display at present."}
{"text": "Rabies is an inevitably fatal but preventable disease. In developing countries rabies remains a significant endemic disease burden. World-wide approximately 55,000 people die of rabies each year (WHO Expert Consultation on Rabies, 2004). Rabies is preventable with proper early post-exposure treatment. Currently, post-exposure prophylaxis includes thorough wound-washing with soap and water followed by administration of vaccine and anti-rabies virus immunoglobulin (RIG) of human or equine origin. RIG administered shortly after exposure at the wound site provides passive immunity which neutralizes rabies virus and prevents its spread until the patient's immune response following vaccination is elicited. Deaths due to post-exposure prophylaxis failure are most commonly attributed to deviations from the recommended regimen such as late initiation of post-exposure prophylaxis or no administration of RIG (Wilde, Vaccine, 25:7605-7609, 2007). In developing countries, availability of RIG is extremely low with only 1-2% of post-exposure prophylaxis being performed using RIG (Sudarshan et al., Int J Infect Dis, 11:29-35, 2007; WHO Consultation on a Monoclonal Antibody Cocktail for Rabies Post Exposure Treatment, 2002). In the United States, only human derived RIG is administered due to the risk of anaphylactic shock from exposure to equine immunoglobulins, but the concern of blood-born pathogen transmission in human RIG remains.\nHuman monoclonal antibodies (mAbs) that neutralize rabies virus have long been recognized as an alternative to overcome the limitations of RIG (Dietzschold et al., J Virol, 64:3087-3090, 1990). Adequate supplies of cell-cultured human mAbs could be produced in a cost-effective manner (Prosniak et al., J Infect Dis, 188:53-56, 2003). In addition, the use of human mAbs reduces the likelihood of an adverse immune response (Weiner, J Immunother, 29:1-9, 2006) and has been shown to be as effective as RIG in preventing rabies in animals (de Kruif et al., Annu Rev Med, 58:359-368, 2007). When a cocktail of two rabies-neutralizing, human mAbs was given with a rabies vaccine to animals experimentally infected with rabies, dose-dependent survival was observed, and all animals receiving the highest dose survived (Goudsmit et al., J Infect Dis, 193:796-801, 2006). The same cocktail was recently shown to be safe to administer to healthy humans in two phase-one clinical trials (Bakker et al., Vaccine, 26:5922-5927, 2008). However, selection of human monoclonal antibodies to include in such a cocktail has some limitations. The diversity of mAbs produced depends significantly on the diversity of viral antigens used to immunize human donors.\nRabies virus is a member (genotype 1) of the genus Lyssavirus. This genus also includes rabies-like viruses (genotypes 2-7) which can cause rabies disease in humans (Bourhy et al., Virology, 194:70-81, 1993; Gould et al., Virus Res, 54:165-187 1998). These viruses have a non-segmented, negative-sense, single-stranded RNA genome that encodes five proteins (Sokol et al., Virology, 38:651-665, 1969; Sokol et al., J Virol, 7:241-249, 1971). In the mature virion, RNA-dependent RNA polymerase, phosphoprotein, and nucleocapsid protein are associated with the genomic RNA while matrix protein and glycoprotein (G protein) surround it (Wiktor et al., J Immunol, 110:269-276, 1973). Trimeric G protein “spikes” coat the surface of the virion and as the only surface exposed protein, is responsible for attachment and entry into host cells. This also makes G protein the primary antigen for induction of virus-neutralizing antibodies, and G protein-specific mAbs are included in the cocktails currently being developed for post-exposure prophylaxis (Dietzschold et al., J Virol, 64:3087-3090, 1990; Kramer et al., Eur J Immunol, 35:2131-2145, 2005; Wunner et al., “Rabies Virus” in Rabies (Second Edition),” pp. 23-68, Academic Press, Oxford, 2007)."}
{"text": "1. Field of the Invention\nThe present invention relates to an inverter controller for controlling a pulse-width modulated inverter. More particularly, the present invention pertains to a controller for controlling inverter switching elements which involves no fear of a short across the power supply during the commutation of a voltage-controlled inverter having a relatively small impedance when viewing the power supply side from the load, and which enables the control accuracy and response to be improved.\n2. Description of the Related Art\nIn pulse-width modulated inverter controllers, an ON signal which is supplied to semiconductor switching elements is generally provided with a dead time of several tens of .mu.sec for the purpose of preventing the occurrence of a short across the power supply during the commutation.\nThis type of conventional inverter controller will be explained below with reference to FIGS. 10 to 13. FIG. 10 shows the circuit configuration of a conventional inverter controller, in which a single-phase inverter is exemplarily illustrated for simplification of the explanation.\nAn inverter 1 is constituted by semiconductor switching elements. More specifically, the inverter 1 is formed by series-connecting an anti-parallel circuit of a transistor 9 and a diode 11 and another anti-parallel circuit of a transistor 10 and a diode 12, these anti-parallel circuits being respectively connected to the positive side of a DC power supply 15 and the negative side of a DC power supply 16. One end of a single-phase load 13 is connected to the node of series connection between the transistor 9 and the transistor 10 which defines the output terminal of the inverter 1. The other end of the load 13 is grounded through a current detector 14. The node of series connection between the DC power supplies 15 and 16 is similarly grounded.\nA current control amplifier 18 is supplied with, as inputs, an output current command value output from an output current command value calculating section 17 and a detected output current value output from the current detector 14, so as to calculate and output an output potential command for the inverter 1.\nA delay circuit 4 is supplied with, as an input, the output potential command from the current control amplifier 18 and outputs a signal obtained by delaying the output potential command.\nAn AND circuit 6 is supplied with, as inputs, the output potential command and the delayed output potential command. The AND circuit 6 ANDs these signals to form an ON/OFF signal for the transistor 9. Another delay circuit 5 is supplied with the output potential command which is output from the current control amplifier 18 and inverted by a NOT circuit 3. The delay circuit 5 delays this inverted output potential command and outputs the delayed command to an AND circuit 7.\nThe AND circuit 7 is supplied with, as inputs, the inverted output potential command output from the NOT circuit 3 and the signal obtained by delaying the inverted output potential command through the delay circuit 5. The AND circuit 7 ANDs these signals and outputs the result of the ANDing as an ON/OFF signal for the transistor 10.\nA driver 8 supplies the base of the transistor 9 with a signal processed on the basis of the ON/OFF signal for the transistor 9 output from the AND circuit 6 and also supplies the base of the transistor 10 with a signal processed on the basis of the ON/OFF signal for the transistor 10 output from the AND circuit 7.\nFIG. 11 is a timing chart of various signals. The output of the delay circuit 4 is a signal which has a predetermined delay .DELTA.T with respect to the output potential command. In consequence, the output of the AND circuit 6, that is, the ON/OFF signal for the transistor 9, rises when a predetermined period of time .DELTA.T has elapsed after the rise of the output potential command. On the other hand, the output of the delay circuit 5 is a signal which has a predetermined delay .DELTA.T with respect to the inverted output potential command. In consequence, the output of the AND circuit 7, that is, the ON/OFF signal for the transistor 10 rises when a predetermined period of time .DELTA.T has elapsed after the fall of the output potential command. As a result, dead time .DELTA.T can be provided between the ON signal for the transistor 9 and the ON signal for the transistor 10. This dead time is provided for the purpose of preventing the occurrence of a short across of the power supply due to a possible delay in operation of the transistors when turned OFF and is generally set at several tens of .mu.sec. The operation of the conventional inverter controller involving such dead time is shown in FIGS. 12 and 13. Referring first to FIG. 12, an ON signal is supplied to the transistor 9 in the mode a, and current is flowing into the load 13 through the transistor 9. When, in this state, an OFF signal is supplied to the transistor 9 (the mode b), the current flowing through the transistor 9 decreases, and the current flowing through the diode 12 gradually increases. In this state, the transistor 9 has not yet been completely turned OFF. Therefore, if an ON signal is supplied to the transistor 10 in this state, the power supply is shorted. For this reason, it is necessary to supply an ON signal to the transistor 10 after the elapse of time which is sufficient for the transistor 9 to turn OFF completely, as shown in the mode c. The period of time during the mode b provided when the operation mode is shifted from the mode a to the mode c is referred to as \"dead time\".\nDead time is also provided when the operation mode is shifted from the mode c in which the transistor 10 is ON to the mode a in which the transistor 9 is ON. More specifically, the mode d is provided between the modes c and a.\nThe inverter operates in the above-described modes a to d when the output current is positive. When the output current is negative also, dead time such as modes f and h shown in FIG. 13 is similarly provided.\nThe dead time is set such to be sufficiently longer than a possible delay in operation of the transistors in order to prevent the occurrence of a short across the power supply and ensure the safety. However, the dead time constitutes an error with respect to the output potential command and therefore leads to lowering in the degree of accuracy in the current control. In addition, since the rise of the ON signal is delayed with respect to the output potential command, the response and stability are deteriorated.\nTo reduce the losses and noise in the load due to harmonic components, it is necessary to increase the switching frequency of the inverter so that the current ripple is reduced. However, since the above-described problems arise every switching operation of the inverter, such problems become even more conspicuous as the switching frequency becomes higher. Accordingly, the upper limit of switching frequency cannot be set at a high value, and the current ripple cannot be sufficiently reduced.\nTo solve the above-described problems, a method has already been tried in which the output voltage is fed back to detect an error portion in the output voltage due to the dead time, and the current control is corrected on the basis of this error portion (see, e.g., the specification of Japanese Patent Laid-Open No. 123478/1984).\nThis method suffers, however, from the following problems. Namely, the current control amplifier 18 shown in FIG. 10 needs to detect the error portion in the output voltage due to the dead time and carry out calculation for correcting the current control on the basis of the detected error. In addition, it is necessary to detect the output voltage at high speed and with a high degree of accuracy. Consequently, the current control method is complicated, and the controller for carrying out the method is costly. Further, since the switching operation of the inverter takes place with a delay corresponding to the dead time with respect to the output potential command at all times, the response and stability cannot be improved. This conventional method further involves a delay by the calculation of an error in the output voltage due to the dead time, and a delay by the calculation for correcting the error. It is therefore impossible to apply this method to the highly responsive current control in which the inverter is directly controlled by the instantaneous value of the current.\nIn order to overcome the above-described problems, another method has heretofore been tried in which the direction of the output current is fed back, and when the output current is positive, the supply of the ON signal to the transistor 10 is inhibited, whereas, when the output current is negative, the supply of the ON signal to the transistor 9 is inhibited, thereby carrying out current control without the need to provide dead time. This method, however, needs to feed back the direction of the output current, which means that the arrangement of the controller is complicated. The method is therefore disadvantageous from the economical point of view. In addition, when the detection system for detecting the direction of the output current is subjected to external disturbance such as noise, a serious accident, i.e., a short across the power supply may be caused, and this means that this method is unsatisfactory in terms of the reliability. When the output current is at or near the zero-crossing point, any drift of the detector leads to an erroneous detection of the direction of the output current, so that it becomes impossible to effect stable control of the inverter. For example, when the direction of the output current is positive and, therefore, the supply of the ON signal to the transistor 9 should be allowed, if the direction of the output current is errorneously recognized to be negative due to a drift of the detector and the supply of the ON signal to the transistor 9 is consequently inhibited, it becomes impossible to increase the output current any more."}
{"text": "Antibody-drug conjugates (ADCs)—monoclonal antibodies (mAbs) covalently linked to toxic agents are acknowledged and employed in the art as improved anticancer treatments. The combination of targeting specificity of mAbs with cytotoxic small molecules, allows for discrimination between malignant and normal tissues resulting in fewer toxic side effects exhibited by many conventional chemotherapies.\nCurrently, there exist only three FDA-approved ADCs available for cancer treatment. Two of the ADCs, Zevalin® and Bexxar®, comprise mAbs covalently attached to the radioisotopes Yttrium-90 and Iodine-131, respectively. These murine radiolabeled mAbs are used in the treatment of CD20-expressing B-cell lymphomas. The third FDA-approved ADC, Mylotarg®, consists of a humanized anti-CD33 monoclonal antibody conjugated to the DNA-cleaving enediyne antibiotic, Calicheamicin. To date, this is the only approved immunoconjugate possessing a cytotoxic organic molecule.\nRecently, groups at the National Cancer Institute (NCI) reported on the preparation and evaluation of immunoconjugates composed of HERCEPTIN® (Trastuzumab) and the cytotoxin, Geldanamycin. Mandler, R., et al., Trastuzumab-Geldanamycin Immunoconjugates: Pharmacokinetics, Biodistribution, and Enhanced Antitumor Activity, Cancer Res. 64:1460-1467 (2004). Geldanamycin is a highly cytotoxic ansamycin benzoquinone antibiotic that exerts its cytotoxicity by binding to the protein chaperone heat shock protein 90 (hsp 90). The anticancer potential of Geldanamycin as a single chemotherapeutic agent has been abandoned as a result of its nonselective and severe toxicity. Geldanamycin is unique in that targeting of hsp90 is effective at down-regulating HER2, and thus tumor cells that over express HER2 are especially sensitive to Geldanamycin. Trastuzumab-Geldanamycin was shown to have IC50s 10-200-fold lower than that of unmodified Trastuzumab in antiproliferative assays. In a xenograft murine model, consisting of weekly i.p. doses of 4 mg/kg over four months, animals treated with Trastuzumab-Geldanamycin exhibited 69% tumor regression, whereas those treated with Trastuzumab alone only showed 7% regression. Median survival time was 145 days as opposed to 78 days for Trastuzumab-treated animals. In addition, nearly one-third of the mice remained tumor free two months after treatment had been completed, whereas there were no survivors from the Trastuzumab group.\nNearly all organic anticancer agents are associated with dose-limiting toxicities. A significant need exists to improve drug efficacy while minimizing systemic toxicity while addressing optimization parameters, such as physiological barriers to targeting molecule extravasation and intratumoral penetration, immunonconjugate aggregation, immunogenicity, normal tissue expression of targeted antigens, and inefficient drug release from the carrier. Certain compounds, moreover, require an improved delivery system due to physicochemical properties such as water solubility, cellular uptake, as well as otherwise short half-life in vivo. Since prognosis of long-term, disease-free survival of most cancer patients remains poor, there continues to be a critical unmet need for improved anticancer treatment."}
{"text": "This invention relates generally to monitoring systems and more particularly concerns devices and systems used to monitor bed patients in hospital or other care giving environments.\nIt is well documented that the elderly and post-surgical patients are at a heightened risk of falling. There are many reasons for this but, broadly speaking, these individuals are often afflicted by gait and balance disorders, weakness, dizziness, confusion, visual impairment, and postural hypotension (i.e., a sudden drop in blood pressure that causes dizziness and fainting), all of which are recognized as potential contributors to a fall. Additionally, cognitive and functional impairment, and sedating and psychoactive medications are also well recognized risk factors.\nA fall places the patient at risk of various injuries including sprains, fractures, and broken bonesxe2x80x94injuries which in some cases can be severe enough to eventually lead to a fatality. Of course, those most susceptible to falls are often those in the poorest general health and least likely to recover quickly from their injuries. In addition to the obvious physiological consequences of fall-related injuries, there are also a variety of adverse economic and legal consequences that include the actual cost of treating the victim and, in some cases, caretaker liability issues.\nIn the past, it has been commonplace to treat patients that are prone to falling by limiting their mobility through the use of restraints, the underlying theory being that if the patient is not free to move about, he or she will not be as likely to fall. However, research has shown that restraint-based patient treatment strategies are often more harmful than beneficial and should generally be avoidedxe2x80x94the emphasis today being on the promotion of mobility rather than immobility. Among the more successful mobility-based strategies for fall prevention include interventions to improve patient strength and functional status, reduction of environmental hazards, and staff identification and monitoring of high-risk hospital patients and nursing home residents.\nOf course, monitoring high-risk patients, as effective as that care strategy might appear to be in theory, suffers from the obvious practical disadvantage of requiring additional staff if the monitoring is to be in the form of direct observation. Thus, the trend in patient monitoring has been toward the use of electrical devices to signal changes in a patient\"\"s circumstance to a caregiver who might be located either nearby or remotely at a central monitoring facility, such as a nurse\"\"s station. The obvious advantage of an electronic monitoring arrangement is that it frees the caregiver to pursue other tasks away from the patient. Additionally, when the monitoring is done at a central facility a single nurse can monitor multiple patients which can result in decreased staffing requirements.\nGenerally speaking, electronic monitors work by first sensing an initial status of a patient, and then generating a signal when that status changes, e.g., he or she has sat up in bed, left the bed, risen from a chair, etc., any of which situations could pose a potential cause for concern in the case of an at-risk patient. Electronic bed and chair monitors typically use a pressure sensitive switch in combination with a separate monitor/microprocessor. In a common arrangement, a patient\"\"s weight resting on a pressure sensitive mat (i.e., xe2x80x9csensingxe2x80x9d mat) completes an electrical circuit, thereby signaling the presence of the patient to the microprocessor. When the weight is removed from the pressure sensitive switch, the electrical circuit is interrupted, which fact is sensed by the microprocessor. The software logic that drives the monitor is typically programmed to respond to the now-opened circuit by triggering some sort of alarmxe2x80x94either electronically (e.g., to the nursing station via a conventional nurse call system) or audibly (via a built-in siren). Some examples of devices that operate in this general fashion may be found in U.S. Pat. Nos. 4,484,043, 4,565,910, 5,554,835, and 5,634,760, the disclosures of which are incorporated herein by reference.\nThat being said, patient monitoring systems that rely on sensing mats to detect the presence of a patient in a bed suffer from a variety of drawbacks. For example, the bed monitoring systems currently available in the marketplace feature externally accessible configuration switches that allow the caregiver to reconfigure the device at will and to adjust parameters such as the duration of the alarm, and the time lapse between the sensing of the xe2x80x9cempty bedxe2x80x9d condition and the sounding of an alarm. External switching makes tampering with the system extremely easy and makes it more difficult to establish and maintain a hospital-wide policy with respect to monitor settings.\nA further problem with conventional bed monitoring systems is that they use oscillating transducers in their alarm audio circuits, resulting in single frequency audio alarms. Since bed monitor alarms are frequently employed in environments in which a multiplicity of other problems might also trigger audio alarms, if the single alarm sound provided by the bed monitor happens to be similar to one or more other alarm sounds heard in response to different monitors, confusion and consequential lengthened response times to patient monitor alarms may result.\nThose skilled in the art know that there are many nurse call station configurations and it is to the economic advantage of a manufacturer to be able to accommodate all of them. However, another problem with the present state-of-the-art in bed monitoring systems is that they are typically pre-configured internally at the factory for one particular type of nurse call station. Thus, if the unit is misconfigured when it arrives at an installation, it may be necessary to summon a medical technician to reconfigure it, since internal modifications to the unit are required to adapt it to different call station types. This can result in additional expense and delay in getting the unit correctly configured and into operation. Further, there are many hospitals that use multiple incompatible nurse call system types, each having been separately added as a new building or wing was constructed. The inability to quickly and reliable move electronic monitors between these systems means that the hospital will generally be required to maintain excess inventory of each type of compatible monitor, a result that ultimately adds to the health care costs borne by the consumer/patient.\nStill another failure in known bed monitoring systems is that they do not provide a method of accumulating statistical data relating to the operation of the unit including, for example, the response times of the caregiver to alarm conditions. This sort of information could be very helpful to the maintenance and proper operation of the monitor, and for caregiver quality control purposes.\nIt is, therefore, a primary object of this invention to provide a patient monitor that is microprocessor-based so as to be reconfigurable by the uploading of configuration data to an electronically erasable programmable read only memory accessible by the microprocessor. A further object of this invention is to provide a microprocessor based patient monitor which synthesizes multiple alarm sounds in software for selection by the caregiver. It is also an object of this invention to provide a microprocessor based patient monitor having a nurse call interface allowing interconnection with any nurse call station without modification of the monitor. Yet another object of this invention is to provide a microprocessor based patient monitor having an electrically erasable programmable read only memory accessible by the microprocessor for logging statistical data with respect to the use of the monitor and the response time of the caregiver using the monitor. Another object of this invention is to provide a microprocessor based bed patient monitor which permits the downloading of the logged statistical data to a host microprocessor connected to the system. It is still another object of the instant invention to provide a system for configuration of monitor parameters and for recalling and analyzing statistical data accumulated therein.\nHeretofore, as is well known in the bed monitor arts, there has been a need for an invention to address and solve the above-described problems. Accordingly, it should now be recognized, as was recognized by the present inventor, that there exists, and has existed for some time, a very real need for a electronic patient monitor that would address and solve the above-described problems.\nBefore proceeding to a description of the present invention, however, it should be noted and remembered that the description of the invention which follows, together with the accompanying drawings, should not be construed as limiting the invention to the examples (or preferred embodiments) shown and described. This is so because those skilled in the art to which the invention pertains will be able to devise other forms of this invention within the ambit of the appended claims.\nIn accordance with the invention, a patient monitor is provided in which a processor receiving electronic signals from a sensor indicating the presence on the sensor and absence from the sensor of a patient is combined with an alarm system which includes a loudspeaker driven by a power amplifier which responds to an input signal derived from a programmable volume control to produce an aural alarm. The processor synthesizes at least one and preferably multiple alarm sounds under software control, operates the programmable volume control of the alarm system to select the decibel level of the alarm and activates and deactivates the alarm in response to the electronic signals received from the sensor and a user interface. An electrically erasable programmable read-only memory accessible by the processor stores a plurality of alarm sounds for selection by the processor for synthesis of the selected alarm sound. In addition, the electrically erasable programmable read-only memory stores multiple decibel levels for selection by the processor of the desired decibel level of the alarm sound. In the preferred embodiment, the patient monitor will be used to sense the presence of patient who is lying in a bed, however, it should be noted and remembered this monitor could also be used in other sorts of applications, including with chair and toilet monitors.\nPreferably, the electrically erasable programmable read-only memory also permits storage of a plurality of options for the delay time between initiation of the absence of a patient from the sensor and the activation of the alarm by the processor. Furthermore, the monitor is preferably provided with an external switch connected to the processor for caregiver selection of the delay time from the plurality of delay time options.\nIt is also preferred that the electrically erasable programmable read-only memory log usage data with respect to the monitor including the total hours of use of the monitor, the total time of alarms sounded by the monitor, the total number of alarms sounded by the monitor and the response time between the most recent sounding of an alarm and a subsequent operation of the monitor by the caregiver. The monitor will include a port for downloading the log usage data to a host computer.\nThe monitor also includes a nurse call interface having a relay which is energized when the power amplifier is de-energized and which has a normally opened contact, a normally closed contact and a common contact for interconnecting the monitor to a nurse call system to one of the normally opened and normally closed contacts so that the monitor requires no modification to accommodate the type of nurse call station with which the monitor is used.\nAccording to still another aspect of the instant invention, there is provided a bed monitor/computer system which allows easy on-site configuration of a monitor to work with different nurses stations. In more particular, the monitor of the instant invention is designed to be reconfigured through the use of a host computer, which obviates the need for internal modifications of monitor parameters through the use of dip switches, rotary dials, etc., which are commonly used in the industry. In the preferred embodiment, a standard computer interface, such as serial interface, is provided as a means for communication between the monitor and a separate host computer. This allows the unit to be readily reprogrammed without risking the exposure of the internal electronic components to the environment.\nAccording to still a further aspect of the instant invention, there is provided a software system for providing the monitor with new programming instructions or a new xe2x80x9cpersonalityxe2x80x9d which will enable it to operate with potentially any plug-compatible nurse call station. In the preferred embodiment, the internal operating logic and various parameters which change the operation of the device to match a particular nurse call station are preferably stored in nonvolatile flash-type RAM which is RAM that can be modified on demand through the use of a host computer-to-patient monitor transfer. One obvious advantage of this arrangement is that it eliminates the many problems associated with mechanical configuration switches, such as dip switches and rotary dials, while providing an easy, inexpensive, and reliable way of upgrading or otherwise modifying the functionality of a monitor while it is in the field.\nThe foregoing has outlined in broad terms the more important features of the invention disclosed herein so that the detailed description that follows may be more clearly understood, and so that the contribution of the instant inventor to the art may be better appreciated. The instant invention is not to be limited in its application to the details of the construction and to the arrangements of the components set forth in the following description or illustrated in the drawings. Rather, the invention is capable of other embodiments and of being practiced and carried out in various other ways not specifically enumerated herein. Additionally, the disclosure that follows is intended to cover all alternatives, modifications and equivalents as may be included within the spirit and scope of the invention as defined by the appended claims. Further, it should be understood that the phraseology and terminology employed herein are for the purpose of description and should not be regarded as limiting, unless the specification specifically so limits the invention.\nWhile the instant invention will be described in connection with a preferred embodiment, it will be understood that it is not intended to limit the invention to that embodiment. On the contrary, it is intended to cover all alternatives, modifications and equivalents as may be included within the spirit and scope of the invention as defined by the appended claims."}
{"text": "1. Field of the Invention\nThe invention concerns an illumination system for wavelengths ≦100 nm, wherein the illumination system has an object plane and a field plane. In illumination systems ≦100 nm, the problem exists that light sources of such illumination systems emit radiation that can lead to an undesired exposure of the light-sensitive object in the wafer plane. Furthermore optical components of the exposure system, such as, for example, the multilayer mirror, can be heated up in this way.\nIn order to filter out the undesired radiation, transmission filters are used in illumination systems for wavelengths ≦100 nm. Such filters have the disadvantage of high light losses. In addition, they can be disrupted very easily by heat stress.\nThe object of the invention is to provide an illumination system for wavelengths ≦100 nm, particularly in the EUV range, in which the above-named disadvantages can be avoided.\nAccording to the invention, this object is solved by an illumination system that has at least one grating element and at least one physical diaphragm in a diaphragm plane. The physical diaphragm is situated in the beam path from the object plane to the field plane after the grating element.\n2. Description of the Prior Art\nGrating elements, for example, reflection gratings, particularly echelette gratings, which are also known as blazed gratings, have been known for a long time from monochromator construction for synchrotron radiation sources. For these elements good experiences, particularly with very high fluxes, were made.\nWith respect to the use of diffraction gratings in monochromators, reference is made to the following publications, whose disclosure content is incorporated to the full extent in the present Application:                H. Petersen, C. Jung, C. Hellwig, W. B. Peatman, W. Gudat: “Review of plane grating focusing for soft x-ray monochromators”, Rev. Sci. Instrum. 66(1), January 1996.        M. V. R. K. Murty: “Use of convergent and divergent illumination with plane gratings”, Journal of the Optical Society of America, Vol. 52, No. 7, July 1962, pp. 768–773.        T. Oshio, E. Ishiguro, R. Iwanaga: “A theory of new astigmatism and coma-free spectrometer”, Nuclear Instruments and Methods 208 (1993) 297–301.        "}
{"text": "The present invention relates to a slurry pump, and more particularly, to a piston assembly or head having an improved self-flushing feature."}
{"text": "The present invention relates to apparatus and methods for providing electrical continuity between two objects, and more particularly to an array of solderless connectors for use with a land grid array integrated circuit package.\nLand grid array (LGA) connector assemblies are commonly used with integrated circuit (IC) packages, such as in applications which do not require soldering of the pins of the LGA connector assembly to either the IC package or a corresponding circuit board. As one example, an LGA connector assembly can be used to temporarily place an LGA package in electrical communication with a circuit card during test, emulation, and debug procedures. As another example, the LGA socket assembly can be used for upgrades and replacements of LGA packages onto circuit boards.\nThe present invention incorporates a variety of novel and unobvious features which are improvements over currently existing LGA socket assemblies.\nOne aspect of the present invention includes an apparatus for providing electrical continuity between two objects. The apparatus includes a body with a top surface and a bottom surface, the body defining a plurality of pin receptacles, each receptacle including a guiding slot within the body between the top and bottom surfaces. The apparatus includes a plurality of pins, each one of the pins being located within a different one of the plurality of receptacles, each pin including a centerbody with two edges, a first member extending from the centerbody, a first cantilever beam extending from the centerbody, and a second cantilever beam extending from the centerbody. The first member of each one of the plurality of pins cooperates with the guiding slot of the corresponding receptacle to guide the pin within the receptacle, each pin being freely moveable within the corresponding receptacle.\nAnother aspect of the present invention includes an apparatus for providing electrical continuity between two objects. The apparatus includes a body with a top surface and a bottom surface, the body defining a plurality of pin receptacles, each receptacle including an aperture. The apparatus includes a plurality of pins, each one of the pins being loose within a different one of the plurality of receptacles, each pin including a centerbody, a first cantilever beam extending from of the centerbody at an acute angle relative to the centerbody, and a second cantilever beam extending from the centerbody at an acute angle relative to the centerbody. The first cantilever beam includes a free end that extends over an adjacent one of the pins.\nAnother aspect of the present invention includes an apparatus for providing electrical continuity between two objects. The apparatus includes a body with a top surface and a bottom surface, the body defining a plurality of pin receptacles, each receptacle including an aperture and a guiding slot within the body. The apparatus includes a plurality of pins located within the plurality of receptacles, each pin including a planar centerbody, a first member extending from the centerbody and cooperating with the guiding slot to loosely locate each pin within a corresponding receptacle, and a first cantilever beam extending from the centerbody. The centerbody includes a projection extending from a surface of the centerbody, the projection cooperating with the receptacle to limit sliding motion of said pin within the receptacle.\nThese and other aspects of the present invention will be apparent from the claims, drawings, and the description of the preferred embodiment to follow."}
{"text": "A number of patents and publications are cited herein in order to more fully describe and disclose the invention and the state of the art to which the invention pertains. Each of these references is incorporated herein by reference in its entirety into the present disclosure, to the same extent as if each individual reference was specifically and individually indicated to be incorporated by reference.\nThroughout this specification, including the claims which follow, unless the context requires otherwise, the word “comprise,” and variations such as “comprises” and “comprising,” will be understood to imply the inclusion of a stated integer or step or group of integers or steps but not the exclusion of any other integer or step or group of integers or steps.\nIt must be noted that, as used in the specification and the appended claims, the singular forms “a,” “an,” and “the” include plural referents unless the context clearly dictates otherwise. Thus, for example, reference to “a pharmaceutical carrier” includes mixtures of two or more such carriers, and the like.\nRanges are often expressed herein as from “about” one particular value, and/or to “about” another particular value. When such a range is expressed, another embodiment includes from the one particular value and/or to the other particular value. Similarly, when values are expressed as approximations, by the use of the antecedent “about,” it will be understood that the particular value forms another embodiment.\nThis disclosure includes information that may be useful in understanding the present invention. It is not an admission that any of the information provided herein is prior art or relevant to the presently claimed invention, or that any publication specifically or implicitly referenced is prior art.\nRAF, Proliferative Disorders, and Cancer\nMutations in genes that directly or indirectly control cell growth and differentiation are generally considered to be the main cause of cancer. Malignant tumors develop through a series of stepwise, progressive changes that lead to the loss of growth control characteristic of cancer cells, i.e., continuous unregulated proliferation, the ability to invade surrounding tissues, and the ability to metastasize to different organ sites. Carefully controlled in vitro studies have helped define the factors that characterize the growth of normal and neoplastic cells and have led to the identification of specific proteins that control cell growth and differentiation.\nRAF is key downstream target for the ras GTPase and mediates the activation of the MAP kinase cascade consisting of raf-MEK-ERK. Activated ERK is a kinase that subsequently targets a number of proteins responsible for mediating, amongst other things, the growth, survival and transcriptional functions of the pathway. These include the transcription factors ELK1, C-JUN, the Ets family (including Ets 1, 2, and 7), and the FOS family. The ras-raf-MEK-ERK signal transduction pathway is activated in response to many cell stimuli including growth factors such as EGF, PDGF, KGF etc. Because the pathway is a major target for growth factor action, the activity of raf-MEK-ERK has been found to be upregulated in many factor dependent tumours. The observation that about 20% of all tumours have undergone an activating mutation in one of the ras proteins indicates that the pathway is more broadly important in tumorigenesis. There is growing evidence that activating mutations in other components of the pathway also occur in human tumours. This is true for RAF.\nThe RAF oncogene family includes three highly conserved genes termed A-RAF, B-RAF and C-RAF (also called Raf-1). RAF genes encode protein kinases that are thought to play important regulatory roles in signal transduction processes that regulate cell proliferation. RAF genes code for highly conserved serine-threonine-specific protein kinases, which are recruited to the plasma membrane following direct binding to the Ras small Guanine-nucleotide binding proteins and this is the initiating event in RAF activation. RAF proteins are part of a signal transduction pathway believed to consist of receptor tyrosine kinases, p21 Ras, RAF protein kinases, Mek1 (ERK activator or MAPKK) kinases and ERK (MAPK) kinases, which ultimately phosphorylate several cellular substrates, including transcription factors. Signaling through this pathway can mediate differentiation, proliferation or oncogenic transformation in different cellular contexts. Thus, RAF kinases are believed to play a fundamental role in the normal cellular signal transduction pathway, coupling a multitude of growth factors to their net effect, cellular proliferation. Because RAF proteins are direct downstream effectors of ras protein function, therapies directed against RAF kinases are believed to be useful in treatment of ras-dependent tumors.\nThe RAF kinases are differentially regulated and expressed; C-RAF is the most thoroughly characterized and is expressed in all organs and in all cell lines that have been examined. A-RAF and B-RAF also appear to be ubiquitous, but are most highly expressed in urogenital and brain tissues, respectively. Because B-RAF is highly expressed in neural tissues it was once thought to be limited to these tissues but it has since been found to be more widely expressed. Although all RAF proteins can bind to active Ras, B-raf is most strongly activated by oncogenic Ras, and may be the primary target of oncogenic Ras in transformed cells.\nRecent evidence indicates that mutational activation of B-RAF is found in a number of different tumours including more than 65% of malignant melanomas, more than 10% of colorectal cancers (Davies, H., et al., 2002, Nature, Vol. 417, pp. 949-954; Rajagopalan, H. et al., 2002, Nature, Vol. 418, p. 934), ovarian cancers (Singer, G., et al., 2003, J. Natl. Cancer Inst., Vol. 95, pp. 484-486) and papillary thyroid cancers (Brose, M., et al., 2002, Cancer Res., Vol. 62, pp. 6997-7000; Cohen, Y., et al., 2003, Invest. Opthalmol. Vis. Sci., Vol. 44, pp. 2876-2878). A range of different B-RAF mutations have been identified in different tumours with the most common being a V600E mutation in the so-called activation loop of the kinase domain (Davies, H., et al., 2002, Nature, Vol. 417, pp. 949-954).\nOther mutations of B-RAF found associated with human cancers may not necessarily activate B-RAF directly but do upregulate the activity of the ras-raf-MEK-ERK pathway by mechanisms which are not fully understood but may involve cross-talk with other RAF isoforms, such as A-RAF (Wan, P., et al., 2004, Cell, Vol. 116, pp. 855-867). In such cases, inhibition of RAF activity would remain a beneficial aim in cancer treatment.\nIn addition to link between B-RAF and certain cancers, there is a significant amount of evidence to indicate a more broad inhibition of RAF activity could be beneficial as an antitumour therapy. Blocking the pathway at the level of B-RAF would be effective at counteracting the upregulation of this pathway caused by tumourigenic ras mutations and also in tumours responding to growth factor action via this pathway. Genetic evidence in Drosophila and C. elegans indicates that RAF homologues are essential for ras dependent actions on differentiation (Dickson, B., et al., 1993, Nature, Vol. 360, pp. 600-603). Introduction of constitutively active MEK into NIH3T3 cells can have a transforming action whilst expression of dominant negative MEK proteins can suppress the tumourigenicity of ras transformed cell lines (Mansour, S. J., et al., 1994, Science, Vol. 265, pp. 966-970; Cowely, S., et al., 1994, Cell, Vol. 77, pp. 841-852). Expression of a dominant negative raf protein has also been found to inhibit ras dependent signalling as has suppression of raf expression using an antisense oligonucleotide construct (Koch, W., et al., 1991, Nature, Vol. 349, pp. 426-428; Bruder, T. T., et al., 1992, Genes and Development, Vol. 6, pp. 545-556).\nThis and other evidence suggests that inhibition of RAF (e.g., B-RAF) activity would be beneficial in the treatment of cancer, and that inhibition of RAF (e.g., B-RAF) activity could be particularly beneficial in those cancers containing a constitutively activated B-raf mutation.\nThe raf-MEK-ERK pathway functions downstream of many receptors and stimuli indicating a broad role in regulation of cell function. For this reason inhibitors of RAF may find utility in other disease conditions that are associated with upregulation of signalling via this pathway. The raf-MEK-ERK pathway is also an important component of the normal response of non-transformed cells to growth factor action. Therefore inhibitors of RAF may be of use in diseases where there is inappropriate or excessive proliferation of normal tissues. These include, but are not limited to glomerulonephritis and psoriasis.\nThe cellular signalling pathway of which RAF is a part has also been implicated in inflammatory disorders characterized by T-cell proliferation (T-cell activation and growth), such as tissue graft rejection, endotoxin shock, and glomerular nephritis.\nRAF (e.g., B-RAF) has been shown to be a valid therapeutic target in hyperproliferative disorders such as cancer. Activated versions of RAF (e.g., B-RAF) are able to transform mammalian cells, allowing them to take on the characteristics of cancer cells and the growth of these cells becomes dependent on the mutant RAF (e.g., B-RAF) protein. Inhibition of RAF (e.g., B-RAF) activity in human cancer cell lines that express the mutant forms of RAF (e.g., B-RAF) blocks their growth and ultimately induces their death.\nAngiogenesis\nChronic proliferative diseases are often accompanied by profound angiogenesis, which can contribute to or maintain an inflammatory and/or proliferative state, or which leads to tissue destruction through the invasive proliferation of blood vessels. (Folkman, 1997, EXS, Vol. 79, pp. 1-81; Folkman, 1995, Nature Medicine, Vol. 1, pp. 27-31; Folkman and Shing, 1992, J. Biol. Chem., Vol. 267, p. 10931.)\nAngiogenesis is generally used to describe the development of new or replacement blood vessels, or neovascularisation. It is a necessary and physiological normal process by which the vasculature is established in the embryo. Angiogenesis does not occur, in general, in most normal adult tissues, exceptions being sites of ovulation, menses and wound healing. Many diseases, however, are characterized by persistent and unregulated angiogenesis. For instance, in arthritis, new capillary blood vessels invade the joint and destroy cartilage (Colville-Nash and Scott, 1992, Ann. Rhum. Dis., Vol. 51, p. 919). In diabetes (and in many different eye diseases), new vessels invade the macula or retina or other ocular structures, and may cause blindness (Brooks et al., 1994, Cell, Vol. 79, p. 1157). The process of atherosclerosis has been linked to angiogenesis (Kahlon et al., 1992, Can. J. Cardiol., Vol. 8, p. 60). Tumor growth and metastasis have been found to be angiogenesis-dependent (Folkman, 1992, Cancer Biol., Vol. 3, p. 65; Denekamp, 1993, Br. J. Rad., Vol. 66, p. 181; Fidler and Ellis, 1994, Cell, Vol. 79, p. 185).\nThe recognition of the involvement of angiogenesis in major diseases has been accompanied by research to identify and develop inhibitors of angiogenesis. These inhibitors are generally classified in response to discrete targets in the angiogenesis cascade, such as activation of endothelial cells by an angiogenic signal; synthesis and release of degradative enzymes; endothelial cell migration; proliferation of endothelial cells; and formation of capillary tubules. Therefore, angiogenesis occurs in many stages and attempts are underway to discover and develop compounds that work to block angiogenesis at these various stages.\nThere are publications that teach that inhibitors of angiogenesis, working by diverse mechanisms, are beneficial in diseases such as cancer and metastasis (O'Reilly et al., 1994, Cell, Vol. 79, p. 315; Ingber et al., 1990, Nature, Vol. 348, p. 555), ocular diseases (Friedlander et al., 1995, Science, Vol. 270, p. 1500), arthritis (Peacock et al., 1992, J. Exp. Med., Vol. 175, p. 1135; Peacock et al., 1995, Cell. Immun., Vol. 160, p. 178) and hemangioma (Taraboletti et al., 1995, J. Natl. Cancer Inst., Vol. 87, p. 293).\nRTKs\nReceptor tyrosine kinases (RTKs) are important in the transmission of biochemical signals across the plasma membrane of cells. These transmembrane molecules characteristically consist of an extracellular ligand-binding domain connected through a segment in the plasma membrane to an intracellular tyrosine kinase domain. Binding of ligand to the receptor results in stimulation of the receptor-associated tyrosine kinase activity that leads to phosphorylation of tyrosine residues on both the receptor and other intracellular proteins, leading to a variety of cellular responses. To date, at least nineteen distinct RTK subfamilies, defined by amino acid sequence homology, have been identified.\nFGFR\nThe fibroblast growth factor (FGF) family of signaling polypeptides regulates a diverse array of physiologic functions including mitogenesis, wound healing, cell differentiation and angiogenesis, and development. Both normal and malignant cell growth as well as proliferation are affected by changes in local concentration of these extracellular signaling molecules, which act as autocrine as well as paracrine factors. Autocrine FGF signaling may be particularly important in the progression of steroid hormone-dependent cancers and to a hormone independentstate (Powers et al., 2000, Endocr. Relat. Cancer, Vol. 7, pp. 165-197).\nFGFs and their receptors are expressed at increased levels in several tissues and cell lines and overexpression is believed to contribute to the malignant phenotype. Furthermore, a number of oncogenes are homologues of genes encoding growth factor receptors, and there is a potential for aberrant activation of FGF-dependent signaling in human pancreatic cancer (Ozawa et al., 2001, Teratoq. Carcinoq. Mutaqen., Vol. 21, pp. 27-44).\nThe two prototypic members are acidic fibroblast growth factor (aFGF or FGF1) and basic fibroblast growth factors (bFGF or FGF2), and to date, at least twenty distinct FGF family members have been identified. The cellular response to FGFs is transmitted via four types of high affinity transmembrane tyrosine-kinase fibroblast growth factor receptors numbered 1 to 4 (FGFR-1 to FGFR-4). Upon ligand binding, the receptors dimerize and auto- or trans-phosphorylate specific cytoplasmic tyrosine residues to transmit an intracellular signal that ultimately reaches nuclear transcription factor effectors.\nDisruption of the FGFR-1 pathway should affect tumor cell proliferation since this kinase is activated in many tumor types in addition to proliferating endothelial cells. The over-expression and activation of FGFR-1 in tumor-associated vasculature has suggested a role for these molecules in tumor angiogenesis.\nFGFR-2 has high affinity for the acidic and/or basic fibroblast growth factors, as well as the keratinocyte growth factor ligands. FGFR-2 also propagates the potent osteogenic effects of FGFs during osteoblast growth and differentiation. Mutations in FGFR-2, leading to complex functional alterations, were shown to induce abnormal ossification of cranial sutures (craniosynostosis), implying a major role of FGFR signaling in intramembranous bone formation. For example, in Apert (AP) syndrome, characterized by premature cranial suture ossification, most cases are associated with point mutations engendering gain-of-function in FGFR-2 (Lemonnier et al., 2001, J. Bone Miner. Res., Vol. 16, pp. 832-845).\nSeveral severe abnormalities in human skeletal development, including Apert, Crouzon, Jackson-Weiss, Beare-Stevenson cutis gyrata, and Pfeiffer syndromes are associated with the occurrence of mutations in FGFR-2. Most, if not all, cases of Pfeiffer Syndrome (PS) are also caused by de novo mutation of the FGFR-2 gene (Meyers et al., 1996, Am. J. Hum. Genet., Vol. 58, pp. 491-498; Plomp et al., 1998, Am. J. Med. Genet., Vol. 75, 245-251), and it was recently shown that mutations in FGFR-2 break one of the cardinal rules governing ligand specificity. Namely, two mutant splice forms of fibroblast growth factor receptor, FGFR2c and FGFR2b, have acquired the ability to bind to and be activated by atypical FGF ligands. This loss of ligand specificity leads to aberrant signaling and suggests that the severe phenotypes of these disease syndromes result from ectopic ligand-dependent activation of FGFR-2 (Yu et al., 2000, Proc. Natl. Acad. Sci. U.S.A., Vol. 97, pp. 14536-14541).\nActivating mutations of the FGFR-3 receptor tyrosine kinase such as chromosomal translocations or point mutations produce deregulated, constitutively active, FGFR-3 receptors which have been involved in multiple myeloma and in bladder and cervix carcinomas (Powers, C. J., et al., 2000, Endocr. Rel. Cancer, Vol. 7, p. 165). Accordingly, FGFR-3 inhibition would be useful in the treatment of multiple myeloma, bladder and cervix carcinomas.\nVEGFR\nVascular endothelial growth factor (VEGF), a polypeptide, is mitogenic for endothelial cells in vitro and stimulates angiogenic responses in vivo. VEGF has also been linked to inappropriate angiogenesis (Pinedo, H. M., et al., 2000, The Oncologist, Vol. 5 (90001), pp. 1-2). VEGFR(s) are protein tyrosine kinases (PTKs). PTKs catalyze the phosphorylation of specific tyrosyl residues in proteins involved in the regulation of cell growth and differentiation. (Wilks, A. F., 1990, Progress in Growth Factor Research, Vol. 2, pp. 97-111; Courtneidge, S. A., 1993, Dev. Supp.l, pp. 57-64; Cooper, J. A., 1994, Semin. Cell Biol., Vol. 5(6), pp. 377-387; Paulson, R. F., 1995, Semin. Immunol., Vol. 7(4), pp. 267-277; Chan, A. C., 1996, Curr. Opin. Immunol., Vol. 8(3), pp. 394-401).\nThree PTK receptors for VEGF have been identified: VEGFR-1 (Flt-1), VEGFR-2 (Flk-1 or KDR), and VEGFR-3 (Flt-4). These receptors are involved in angiogenesis and participate in signal transduction (Mustonen, T., et al., 1995, J. Cell Biol., Vol. 129, pp. 895-898).\nOf particular interest is VEGFR-2, which is a transmembrane receptor PTK expressed primarily in endothelial cells. Activation of VEGFR-2 by VEGF is a critical step in the signal transduction pathway that initiates tumour angiogenesis. VEGF expression may be constitutive to tumour cells and can also be upregulated in response to certain stimuli. One such stimuli is hypoxia, where VEGF expression is upregulated in both tumour and associated host tissues. The VEGF ligand activates VEGFR-2 by binding with itsextracellular VEGF binding site. This leads to receptor dimerization of VEGFRs and autophosphorylation of tyrosine residues at the intracellular kinase domain of VEGFR-2. The kinase domain operates to transfer a phosphate from ATP to the tyrosine residues, thus providing binding sites for signalling proteins downstream of VEGFR-2 leading ultimately to initiation of angiogenesis (McMahon, G., 2000, The Oncologist, Vol. 5(90001), pp. 3-10).\nInhibition at the kinase domain binding site of VEGFR-2 would block phosphorylation of tyrosine residues and serve to disrupt initiation of angiogenesis.\nTIE\nAngiopoieten 1 (Ang1), a ligand for the endothelium-specific receptor tyrosine kinase TIE-2 is a novel angiogenic factor (Davis et al., 1996, Cell, Vol. 87, pp. 1161-1169; Partanen et al., 1992, Mol. Cell Biol., Vol. 12, pp. 1698-1707; U.S. Pat. Nos. 5,521,073; 5,879,672; 5,877,020; and 6,030,831). The acronym TIE represents “tyrosine kinase containing Ig and EGF homology domains”. TIE is used to identify a class of receptor tyrosine kinases, which are exclusively expressed in vascular endothelial cells and early hemopoietic cells. Typically, TIE receptor kinases are characterized by the presence of anEGF-like domain and an immunoglobulin (IG) like domain, which consists of extracellular folding units, stabilized by intra-chain disulfide bonds (Partanen et al., 1999, Curr. Topics Microbiol. Immunol., Vol. 237, pp. 159-172). Unlike VEGF, which functions during the early stages of vascular development, Ang1 and its receptor TIE-2 function in the later stages of vascular development, i.e., during vascular remodelling (remodelling refers to formation of a vascular lumen) and maturation (Yancopoulos et al., 1998, Cell, Vol. 93, pp. 661-664; Peters, K. G., 1998, Circ. Res., Vol. 83(3), pp. 342-343; Suri et al., 1996, Cell, Vol. 87, pp. 1171-1180).\nConsequently, inhibition of TIE-2 would be expected to serve to disrupt remodelling and maturation of new vasculature initiated by angiogenesis thereby disrupting the angiogenic process.\nEph\nThe largest subfamily of receptor tyrosine kinases (RTKs), the Eph family, and their ligands (ephrins), play important roles in physiologic and pathologic vascular processes. Both the Ephs (receptors) and ephrins (ligands) are divided into two groups, A and B subfamilies (Eph Nomenclature Committee, 1997). The binding of ephrin ligands to Eph receptors is dependent on cell-cell interactions. The interactions of ephrins and Ephs have recently been shown to function via bi-directional signalling. The ephrins binding to Eph receptors initiate phosphorylation at specific tyrosine residues in the cytoplasmic domain of the Eph receptors. In response to Eph receptor binding, the ephrin ligand also undergoes tyrosine phosphorylation, so-called ‘reverse’ signalling (Holland, S. J., et al., 1996, Nature, Vol. 383, pp. 722-725; Bruckner et al., 1997, Science, Vol. 275, pp. 1640-1643).\nEph RTKs and their ephrin ligands play important roles in embryonic vascular development. Disruption of specific Eph receptors and ligands (including ephrin-B2) leads to defective vessel remodelling, organisation, and sprouting resulting in embryonic death (Wang, H. U., et al., 1998, Cell, Vol. 93, pp. 741-753; Adams, R. H., et al., 1999, Genes Dev, Vol. 13, pp. 295-306; Gale and Yancopoulos, 1999, Genes Dev, Vol. 13, pp. 1055-1066; Helbling, P. M., et al., 2000, Development, Vol. 127, pp. 269-278). Coordinated expression of the Eph/ephrin system determines the phenotype of embryonic vascular structures: ephrin-B2 is present on arterial endothelial cells (ECs), whereas EphB4 is present on venous ECs (Gale and Yancopoulos, 1999, Genes Dev, Vol. 13, pp. 1055-1066; Shin, D., et al., 2001, Dev Biol, Vol. 230, pp. 139-150). Recently, specific Ephs and ephrins have been implicated in tumour growth and angiogenesis.\nThe Ephs and ephrins have been found to be overexpressed in many human tumours. In particular, the role of EphB2 has been identified in small cell lung carcinoma (Tang, X. X., et al., 1999, Clin Cancer Res, Vol. 5, pp. 455-460), human neuroblastomas (Tang, X. X., et al., 1999, Clin Cancer Res, Vol. 5, pp. 1491-1496) and colorectal cancers (Liu, W., et al., 2004, Brit. J. Canc., Vol. 90, pp. 1620-1626), and higher expression levels of Ephs and ephrins, including EphB2, have been found to correlate with more aggressive and metastatic tumours (Nakamoto, M. and Bergemann, A. D., 2002, Microsc. Res Tech, Vol. 59, pp. 58-67).\nConsequently, inhibition of EphB2 would be expected to serve to disrupt angiogenesis, and in particular in certain tumours where over-expression occurs.\nThe inventors have discovered compounds that, e.g., inhibit RAF (e.g., B-RAF) activity and/or are useful in the treatment of, e.g., proliferative disorders, cancer, etc."}
{"text": "Patient compliance is an important factor to recovery after a surgery. New devices can provide numerous healing benefits but only when the patient follows through with extra tasks. Thus, there is a need in the surgical implant field to create a new and useful system and method for an electrical implant device with increased patient compliance. This invention provides such a new and useful system and method."}
{"text": "Direct sampling receivers convert received signals directly to digital signals for subsequent processing. Direct sampling receivers offer flexibility of use over multiple frequency bands and can improve performance using digital filters that replace less accurate analog counterparts. Speed limitations of analog-to-digital converters (ADCs) implemented within direct sampling receivers limit the performance of ADCs at high sampling frequencies, especially when higher resolution is required.\nComparators are the fundamental building blocks of many types of ADCs and generally dictate the power, performance, and speed of operation of an ADC. Comparators have limitations that lead to significant limitations in the performance and power consumption of ADCs, which negatively impact the overall performance of digital sampling receivers and other devices that include ADCs."}
{"text": "Genetic Engineering of Plants\nThe hurdle of creating successful genetically engineered plants in major crop varieties is now being overcome sequentially on a plant-by-plant basis. The term “genetic engineering” is meant to describe the manipulation of the genome of a plant, typically by the introduction of a foreign gene into a plant, or the modification of the genes of the plant, to increase or decrease the synthesis of gene products in the plant. Typically, genes are introduced into one or more plant cells which can be cultured into whole, sexually competent, viable plants which may be totally transformed or which may be chimeric, having some tissues transformed and some not. These plants can be self-pollinated or cross-pollinated with other plants of the same or compatible species so that the foreign gene or genes carried in the germ line can be bred into agriculturally useful plant varieties.\nCurrent strategies directed toward the genetic engineering of plant lines typically involve two complementary processes. The first process involves the genetic transformation of one or more plant cells of a specifically characterized type. The term “transformation” as used herein means that a foreign gene, typically in the form of a genetic construction, is introduced into the genome of the individual plant cells.\nThis introduction is accomplished through the aid of a vector, which is integrated into the genome of the plant. The second process then involves the regeneration of the transformed plant cells into whole sexually competent plants. Neither the transformation nor regeneration process need to be 100% successful, but must have a reasonable degree of reliability and reproducibility so that a reasonable percentage of the cells can be transformed and regenerated into whole plants.\nEP-A0342926 (the content of which is incorporated herein by reference) discloses a plant (maize) ubiquitin regulatory system comprising a heatshock element (comprising two overlapping consensus heatshock elements), a promoter, a transcription start site, an intron and a translation start site. The heatshock element component of this regulatory system is believed to confer heat inducibility of expression of associated DNA sequences in dicot or monocot cells following permissive levels of heatshock.\nPlant ubiquitin regulatory system refers to the approximately 2 kb nucleotide sequence 5′ to the translation start site of the ubiquitin gene and comprises sequences that direct initiation of transcription, control of expression level, induction of stress genes and enhancement of expression in response to stress. The regulatory system, comprising both promoter (of about 1 kb nucleotide sequence) and regulatory functions, is the DNA sequence providing regulatory control or modulation of gene expression. A structural gene placed under the regulatory control of the plant ubiquitin regulatory system means that a structural gene is positioned such that the regulated expression of the gene is controlled by the sequences comprising the ubiquitin regulatory system.\nPromoters are DNA elements that direct the transcription of RNA in cells. Together with other regulatory elements that specify tissue and temporal specificity of gene expression, promoters control the development of organisms.\nThere has been a concerted effort in identifying and isolating promoters from a wide variety of plants and animals, especially for those promoters demonstrating a high level of constitutive expression and capable of maintaining stable levels of said expression under stress conditions.\nThe present invention is based on modifications of the plant ubiquitin regulatory system."}
{"text": "The invention relates to spectral shift reactor control and more particularly to mechanical means for spectral shift reactor control.\nIn typical nuclear reactors, reactivity control is accomplished by varying the amount of neutron absorbing material (poisons) in the reactor core. Generally, neutron absorbing control rods are utilized to perform this function by varying the number and location of the control rods with respect to the reactor core. In addition to control rods, burnable poisons and poisons dissolved in the reactor coolant can be used to control reactivity.\nIn the conventional designs of pressurized water reactors, an excessive amount of reactivity is designed into the reactor core at start-up so that as the reactivity is depleted over the life of the core the excess reactivity may be employed to lengthen the core life. Since an excessive amount of reactivity is designed into the reactor core at the beginning of core life, neutron absorbing material such as soluble boron must be placed in the core at that time in order to properly control the excess reactivity. Over the core life, as reactivity is consumed, the neutron absorbing material is gradually removed from the reactor core so that the original excess reactivity may be used. While this arrangement provides one means of controlling a nuclear reactor over an extended core life, the neutron absorbing material used during core life absorbs neutrons and removes reactivity from the reactor core that could otherwise be used in a more productive manner such as in plutonium fuel production. The consumption of reactivity in this manner without producing a useful product results in a less efficient depletion of uranium and greater fuel costs than could otherwise be achieved. Therefore, it would be advantageous to be able to extend the life of the reactor core without suppressing excess reactivity with neutron absorbing material thereby providing an extended core life with a significantly lower fuel cost.\nOne such method of producing an extended core life while reducing the amount of neutron absorbing material in the reactor core is by the use of \"Spectral Shift Control\". As is well understood in the art, in one such method the reduction of excess reactivity (and thus neutron absorbing material) is achieved by replacing a large portion of the ordinary reactor coolant water with heavy water. This retards the chain reaction by shifting the neutron spectrum to higher energies and permits the reactor to operate at full power with reduced neutron absorbing material. This shift in the neutron spectrum to a \"hardened\" spectrum also causes more of the U.sup.238 to be converted to plutonium that is eventually used to produce heat. Thus, the shift from a \"soft\" to a \"hard\" spectrum results in more neutrons being consumed by U.sup.238 in a useful manner rather than by poisons. As reactivity is consumed, the heavy water is gradually replaced with ordinary water so that the reactor core reactivity is maintained at a proper level. By the end of core life, essentially all the heavy water has been replaced by oridinary water while the core reactivity has been maintained. Thus, the reactor can be controlled without the use of neutron absorbing material and without the use of excess reactivity at start-up which results in a significant uranium fuel cost savings. The additional plutonium production also reduces the U.sup.235 enrichment requirements. While the use of heavy water as a substitute for ordinary water can be used to effect the \"spectral shift\", the use of heavy water can be an expensive and complicated technology.\nWhile there exist in the prior art numerous ways of controlling a nuclear reactor, what is needed is apparatus for controlling reactor core moderation in a manner that provides for reduced uranium fuel costs and for an extended reactor core life."}
{"text": "The present invention relates to an apparatus for automatically examining and inspecting white resin powder of, for instance, vinyl chloride resins, ABS resins and MBS resins for the presence of foreign substances.\nWhen preparing resin powder or forwarding the powdery product from a manufacturing plant, the product is in general subjected to quality inspection for various predetermined properties. The results are sent back to the manufacturing plant for the improvement of production processes or they are used in the denoration of the quality in order to afford convenience to the destination or the consignee.\nIn case of, for instance, vinyl chloride resin powder, one of the items for quality inspection thereof is to inspect the resin powder for the presence of foreign substances as contaminants. Conventionally, the inspection of the powdery resin product for the presence of foreign substances has been carried out by spreading a constant amount of the vinyl chloride resin powder and searching for the presence of pigmented foreign substances (such as those colored black, brown and/or red) thus to count the number of foreign substances present therein. However, the results widely vary depending on the visual power and the nature of each inspector. Accordingly, this method never provides objectively measured values, requires much labor and time and accordingly, has low efficiency.\nIncidentally, there has recently been used an industrial apparatus for automatic visual inspection which makes use of a video camera. An example of such apparatuses is one which is commercially available from Hajime Sangyo Co., Ltd. under the trade name of FF 4000. This apparatus serves to detect, for instance, defects, stains and foreign substances on the image of a powdery resin sample photographed by the video camera and to specify the positions of the foreign substances on the photograph through blinking modulated bright spots on a monitor. This inspection apparatus permits automatic inspection of a subject for the presence of, for instance, foreign substances and can provide measured values free of any scattering irrespective of inspectors. Moreover, the apparatus does not require much labor and time and ensures efficient inspection.\nHowever, when the aforementioned industrial apparatus for automatic visual inspection which makes use of a video camera is used in the inspection of resin powder easily electrostatically charged such as vinyl chloride resin powder for the presence of colored foreign substances as contaminants thereof, it is often observed that the resin powder to be inspected is not uniformly dropped on a belt conveyor or is unevenly distributed on the belt conveyor due to electrostatic charges generated on the resin powder per se. For this reason, the surface of the spread resin powder layer becomes severely uneven and this leads to a large scattering in the measured values. Consequently, the apparatus does not provide any reproducible result."}
{"text": "The present invention relates generally to the field of electrical measurements, and more particularly to the measurement of nearby electrical and magnetic fields using sensors on a computing device.\nAs provided by the Electrical Safety Foundation International (ESFI), “Contact with Overhead Power Lines” (including direct worker contact and contact through machines, tools, and hand-carried metallic objects) was the largest fatal accident category comprising 44% of all electrical fatalities for the period 1992-2010. “Contact with Wiring, Transformers, or Other Electrical Components” (most common for workers who install, repair, or maintain electrical systems and apparatus in the normal course of their electrical work) was the second-largest fatal injury category with 27% of all fatalities, followed by “Contact with Electric Current of Machine, Tool, Appliance, or Light Fixture” (most common for workers who use electrical tools and apparatus in the normal course of their non-electrical work) with 17%.\nFive occupational groups account for nearly 80% of all fatal electrical accidents. “Construction Trades Workers” represent about 38% of all electrical fatalities. Other occupational groups with high numbers of electrical fatalities are “Installation, Maintenance, and Repair Occupations” (21%), “Other Management Occupations” (4%), and “Agricultural Workers” (2%). The Construction industry experienced 52% of total electrical fatalities."}
{"text": "Computer systems are designed to process a variety of applications, each comprised of software instructions. With the increasing complexity of applications, software instructions have become longer, thus requiring an increased amount of time for the software instructions to be executed.\nSoftware instructions are executed by a microprocessor which is the key working unit of a computer system. The methods that have been used to increase speed in the personal computer have generally centered on maximizing the efficiency with which a single microprocessor can process instructions. Limits are being reached on single microprocessor processing speed. To address this constraint, multiple microprocessors have been combined to operate in parallel within computer systems. Such multiple microprocessor systems, such as symmetrical multiprocessor (i.e., SMP) systems, allocate processing tasks among the multiple parallel microprocessors.\nUpon powering on of an SMP system, each microprocessor is initialized, and a primary microprocessor is selected from among the microprocessors to take charge of bringing up the SMP system. The selection of the primary microprocessor can be done by a software algorithm, or by a hardware locking mechanism. Depending on the number of microprocessors in an SMP system, the time required for initialization increases with the number of microprocessors in the SMP system. The selection of the primary microprocessor can vary upon every time the SMP system is powered on."}
{"text": "The utility of aqueous polymeric dispersions (also called aqueous polymeric emulsions) in the preparation of paints, coatings, adhesives and caulks or sealants is well known. Water based or water borne polymer dispersions are often preferred because of theft relatively low cost, ease of application and relatively low amounts of volatile organic compounds (VOC) contained therein. There is an increasing need for higher solids dispersions which will provide faster setting times for use on high speed production equipment. High solids adhesive bases may also find use as replacements for conventional hot melt (100% solids) material which requires elevated temperatures with consequent expenditures of considerable energy. In addition to the need for high solids dispersion for such applications, it is also essential that the dispersions remain sufficiently low in viscosity that they can be applied using conventional equipment. Considerable effort has been expended to provide aqueous dispersions of polymeric materials which are characterized by a high solids content.\nPressure-sensitive adhesives find use in a wide variety of applications, such as automotive, aerospace, construction and electrical markets, either in the form of tapes or as adhesive coatings on other backings. As the society becomes more aware of the significance of environmental protection, industrial products such as adhesive that is harmful to the ecological environment are gradually being replaced and eliminated.\nCoated surfaces, especially those coated with modern low VOC coatings containing 100 g/L. VOC or less, such as moldings and panels often become visibly discolored after contact with water repeatedly or over extended periods of time. This problem is enhanced when high polymer content coatings are used on the surfaces. Other issues that can occur with painted surfaces are blistering and surfactant leaching, wherein water-soluble components are extracted from the coatings and deposited on the coated surface. Minimizing the amount of water-soluble ingredients has been used to reduce surfactant leaching. However, the water-soluble components are generally substituted with volatile organic compounds (VOCs), resulting in an environmentally undesirable product.\nSurfactants have widely been used as emulsification, dispersion, cleaning, wetting and foaming. Emulsifiers for emulsion polymerization, which are used upon producing polymers by emulsion polymerization, are known not only to take part in polymerization-initiating reactions and polymer-forming reactions but also to affect the mechanical stability, chemical stability, freezing stability, storage stability and the like of the resulting emulsions. Further, they are also known to give significant effects on physical properties of the emulsions, such as particle size, viscosity and foaming potential and, when formed into films, physical properties of the films, such as waterproofness, weatherability, adhesion and heat resistance.\nSurfactants can also be used as reactants, often called reactive surfactants. Surfactants used in polymerization are also called polymerizable and/or copolymerizable surfactants. Among the reactive surfactants, those containing one or more phenyl ether groups as hydrophobic groups have found wide-spread utility for their excellent properties such as emulsifying property, dispersing property, and polymerization-stabilizing property. Thus, surfactants that are based on alkyl phenol ethoxylates (APES) have been widely used in emulsion polymerization. In recent years, however, a concern has arisen about a potential problem that nonyl phenol ethoxylates may show false hormone effects on organisms to disrupt the endocrine system, that is, the so-called endocrine problem has arisen, so that research has also been conducted in efforts to provide replacements for the reactive surfactants containing one or more phenyl ether groups.\nAccordingly, there is a need in the art for an aqueous dispersion and/or emulsion with ultra-high solids but environmentally friendly in preparations and applications."}
{"text": "This invention relates to non-lethal weapons; and more particularly, to a electrical stun gun and electrically conductive liquids."}
{"text": "The present invention relates to tank pressure control apparatus, and particularly to venting apparatus which operates to seal the vapor space in the interior of a fuel tank during refueling and reopens at some point after refueling has been completed to vent the vapor space. More particularly, the present invention relates to a fill limit system for preventing overfilling of a fuel tank.\nFuel pump nozzles are known to include a fill-limiting sensor for shutting off the flow of fuel from the nozzle when the fuel tank is nearly filled. Typically, this fill-limiting sensor is triggered whenever the fuel tank is full and fuel \"backs up\" the filler neck to splash onto or reach the fill-limiting sensor located on the nozzle. Sometimes fuel pump operators overfill a fuel tank inadvertently in a good-faith effort to fill the tank \"completely\" or to purchase a quantity of fuel that can be paid for in cash without causing the operator to receive unwanted coinage in change.\nIt has been observed that fuel pump operators are able to manually override or bypass some fill-limiting sensors on nozzles by continuing to pump fuel after the pump nozzle has automatically shut off several times. This practice has come to be called the \"trickle fill\" method of introducing liquid fuel into a fuel tank. Using this well-known trickle fill method, the fuel pump operator \"clicks\" or squeezes the lever handle on the pump nozzle slowly two or three times in succession after automatic nozzle shutoff has occurred in order to introduce more fuel into the fuel tank without actuating the fill-limiting sensor right away. It will be appreciated that such trickle fill practices can result in overfilling the fuel tank which can effectively reduce the fuel vapor expansion capacity in the vapor space available within the filled fuel tank.\nVehicle fuel systems are known to include valves for venting the vapor space in a fuel tank. See, for example, U.S. Pat. Nos. 4,760,858; 4,991,615; and 5,028,255.\nWhat is needed is a fill limit control system that is able to vent fuel vapor from the vapor space in a fuel tank during early stages of refueling but block introduction of any and all liquid fuel in excess of a maximum volume so as to preserve a minimum volume of the vapor space in the fuel tank once the fuel tank is filled to its maximum capacity with fuel. Ideally, this fill limit control system could be included as one component in a comprehensive Onboard Refueling Vapor Recovery (ORVR) vehicle fuel system. An ORVR system is used to manage fuel vapor recovery during all phases of vehicle use. In addition, a fill limit control system that is adaptable to permit customers some latitude to trickle-fill their fuel tanks without overfilling the fuel tanks would be welcomed by fuel customers and vehicle makers.\nAccording to the present invention, an apparatus is provided for controlling discharge of fuel vapor from a vehicle fuel tank. The apparatus includes a valve housing mounted in a wall of the fuel tank. The valve housing is formed to include a vapor inlet opening communicating with an interior region of the fuel tank, a vapor outlet opening, and a passageway interconnecting the vapor inlet opening and the vapor outlet opening.\nThe apparatus also includes a float valve and a skirt. The float valve is disposed in the valve housing and movable between an open position allowing flow of tank fuel vapor through the vapor outlet opening and a closed position preventing flow of tank fuel vapor through the vapor outlet opening. The skirt is positioned to lie inside the fuel tank and extend around the valve housing and the vapor inlet opening in the valve housing to define an annular channel lying around the valve housing and conducting fuel vapor from the interior region of the fuel tank to the vapor inlet opening.\nThe skirt has a bottom edge lying in spaced-apart relation to the top wall of the fuel tank to locate the vapor inlet opening in the valve housing in a position between the top wall of the fuel tank and the bottom edge. The bottom edge of the skirt is positioned to lie in coextensive relation with a top surface of liquid fuel in the interior region of the fuel tank once the volume of liquid fuel in the fuel tank is equal to a predetermined maximum fuel capacity of the fuel tank.\nIn preferred embodiments, the skirt includes a top wall coupled to the valve housing and an annular side wall appended to the top wall. The annular side wall is positioned to surround the valve housing and define the annular channel therebetween. During refueling, fuel vapor in the fuel tank vapor space will be vented from the fuel tank along a path through the annular channel, vapor inlet opening, passageway, and the vapor outlet opening. Such venting will stop automatically once liquid fuel in the fuel tank rises to a level that is high enough to occlude or close the opening into the annular channel.\nIn effect, the skirt causes the pump nozzle to shut off at the proper time during refueling due to sudden restriction of the exit vapor flow as the rising level of liquid fuel in the fuel tank covers the annular opening into the annular channel at the bottom of the skirt. This allows the fuel pump operator a small amount of trickle fill without materially raising the liquid fuel level in the fuel tank.\nThe annular side wall of the skirt extends downwardly to position the bottom edge in spaced-apart relation to the top wall of the fuel tank. During refueling, the rising liquid fuel will reach the bottom edge of the skirt and close the annular opening into the annular channel before it ever reaches the top wall of the fuel tank. Once the opening into the annular channel is closed, no more fuel vapor can be vented from the vapor space in the fuel tank through the annular channel, vapor inlet opening, passageway, and vapor outlet opening. If a pump operator continues to pump fuel into the fuel tank through the filler neck, the pressure in the vapor space will rise and force some liquid fuel in the fuel tank up through the filler neck to reach and trigger the fill-limiting sensor on the pump nozzle.\nAdvantageously, it is possible to set the maximum fuel capacity of any fuel tank easily by simply mounting a fill-limiting valve assembly in accordance with the present invention so that the bottom edge of the skirt is fixed to lie in coextensive relation with a top surface of liquid fuel in the fuel tank once the volume of liquid fuel in the fuel tank is equal to a predetermined maximum fuel capacity of the fuel tank. It will be understood that use of a bottom edge of the skirt to prevent overfilling of the fuel tank also functions to preserve a minimum volume for the vapor space between the top wall of the fuel tank and the top surface of liquid fuel in the fuel tank.\nIn use, the float valve is expected to rise in the valve housing and close the vapor outlet opening only if a pump operator uses the trickle fill method to try and introduce more liquid fuel into the tank once the automatic shutoff sensor on the pump nozzle has operated to shut off the pump. In such a case, any extra fuel added into the fuel tank will cause the fuel level inside the annular channel to rise, thereby causing the float valve to rise to its closed position. Illustratively, a spring underneath the float valve provides a float-lifting force in addition to buoyancy forces generated by the rising liquid fuel to move the float valve in the annular channel to its closed position.\nAvantageously, a limited amount of trickle fill is permitted by a fill limit valve in accordance with the present invention. In use, after the first automatic pump nozzle shutoff, a pump operator can click the pump nozzle a couple of times to introduce a bit more fuel into the tank. As noted above, this extra fuel displaces fuel already in the tank causing the level of fuel in the annular channel to rise without raising the level of the fuel in the fuel tank outside of the annular channel. In another embodiment, a vertical slot or a small aperture can be formed at the bottom edge of the skirt to allow additional limited trickle fill capability for pump operators without overfilling the tank.\nAdditional features and advantages of the invention will become apparent to those skilled in the art upon consideration of the following detailed description of preferred embodiments exemplifying the best mode of carrying out the invention as presently perceived."}
{"text": "The present invention relates to fuel vapor vent valves employed in vapor emission controls systems in motor vehicle fuel tank installations. The invention particularly relates to float operated vent valves attached to the upper wall of a non-metallic fuel tank. Valves employed for such applications are commonly mounted onto the tank through an access opening formed in the upper wall with portions of the valve extending outwardly over the opening and attached to the outer surface of the tank in a sealing arrangement, as for example, by weldment of compatible non-metallic material to the surface of the tank.\nHeretofore the weldment has been by hot-plate or sonic welding techniques which have proven relatively costly in mass production. Heretofore, attempts to spin weld a fuel vapor vent valve of the present type have resulted in destruction of the valving surfaces by the extreme rotational accelerations and decelerations. Furthermore, proper or accurate orientation of the vent hose connector was not possible with spin welding. However, spin welding has been desired as a more cost effective way of attaching the vent valve to the tank than hot-late or sonic welding."}
{"text": "Hitherto known incontinence shields normally comprise conventional diapers, which are substantially planar and, for example, of rectangular shape. Such diapers, however, are intended for the absorption of urine and excretement, and are thus unsuitable for persons who solely require a urine-collecting incontinence shield or guard. The diaper must have a certain width and the absorbent body must have a given thickness, in order to have satisfactory absorption ability. Naturally, the diaper takes up an equivalent amount of room between the legs of the wearer, causing discomfort to the wearer, in the form of chafing, sores etc. In the case of diapers, there is also a serious risk that excreted body fluids will leak past the edges of the diaper, when the diaper becomes saturated and is compressed between the legs of the wearer during wear.\nIt is known to give the incontinence shield a container-like form, in which the shield comprises an outer layer of liquid-impermeable material and an absorbent material arranged therein. Shields of this kind function quite satisfactorily in the case of men with a penis of normal size, which can be accomodated in the container and embraced thereby. Under certain conditions, however, and particularly when the penis of the wearer is relatively small, there is a serious risk of the penis leaving its position in the container.\nOne particular problem is that of providing a suitable incontinence shield for older men whose penises are retarded and so small that they cannot be accomodated satisfactorily in a hose-like casing.\nThe SE Patent Specification No. 426 206 describes an incontinence shield in which, in an attempt to solve the problem of holding the container-like shield in position also with men whose penis is small or retarded, said shield is arranged to be brought into sealing abutment with the body of the wearer in the area around the root or base of the penis, thereby to try to hold the container more securely.\nThe shield according to this publication has the form of a flat hose-like member having two pairs of mutually opposing edges. The hose-like member has the form of a parallelogram, of which three edges are closed and the fourth edge has been left open. Thus, there is obtained at the opening an acute corner and an obtuse corner. When the shield is worn, the acute corner is turned to face upwards and the obtuse corner is turned to face downwards. As opposed to other known urine-collecting containers in which the opening edge lies at right angles to the long axis of the container, in the shield according to the aforesaid publication the side edges of said opening extend obliquely upwards, along the sides of the root of the penis, it being maintained that in this way the penis is held relatively securely in position, even when the penis of the wearer is relatively small.\nThe aforedescribed incontinence shields, which are solely intended to embrace the penis of the wearer, do not function entirely satisfactorily and are readily able to slip from their intended position.\nThe present invention is based on the concept of using the scrotum or testicle-pouch of the wearer, in addition to his penis, as a means for holding an incontinence shield safely around said penis.\nAccordingly, this invention is characterized in that the shield has a front-piece, a back-piece and a centre-piece (or a connecting line) which connects said front-piece and said back-piece and which is narrower in relation thereto and which serves as the bottom of the container-like shield; in that the shield is arranged to embrace both the penis and the testicle-pouch of the wearer; in that to this end the back-piece is provided with a slit or recess which extends from the end edge of said back-piece; and in that wing-like portions formed on the back-piece by said slit or said recess and located on both sides thereof are intended to be clamped firmly against the body of the wearer, beneath the testicle-pouch thereof. The absorbent body of the shield is suitably soft and compressible, and in addition to functioning as an absorption medium is also intended to impart a given requisite stability to the wing-like portions, to enable said wing-like portions to be compressed and moulded beneath the testicle-pouch of the wearer, so as to hold the shield firmly.\nBecause the shield is intended to embrace both the penis and testicle-pouch of the wearer, the shield will be held more positively, even when the penis of the wearer is small.\nAccording to one suitable embodiment, the inner and outer sides of the shield have portions which extend laterally beyond the absorbent body and there joined together.\nIn order to form a closed, liquid-tight container, the outer edges of said portions on the front-piece of the shield are sealingly joined to corresponding portions on the back-piece thereof.\nThe lateral extension of the absorbent body in the centre-piece should not be too great, since the crutch of the wearer might otherwise become chafed. When seen longitudinally, the size of the front-piece and back-piece of the shield is chosen so as to accomodate both the penis and the testicle-pouch of the wearer. In order to avoid chafing of the crutch of the wearer, the lateral extension of the absorbent body shall be less than one-half the length of the front-piece and back-piece respectively, and suitably greater than one-quarter of said length, in order to provide a satsifactory absorption capacity.\nIn order that the centre-piece of the shield is able to accomodate the testicle-pouch of the wearer, the inner and outer sides of the shield shall have a lateral extension which is substantially from 1.5 to 2.0 times the extension of the absorbent layer.\nIn accordance with one suitable embodiment, a pre-tensioned elastic thread or the like is arranged across the centre-piece, in a manner to gather the sharp edges extending beyond the absorbent layer and formed by the outer and inner sides of the shield.\nThe recess formed in the back-piece of the shield is substantially of U-shape, thereby to obtain close, sealing abutment with the body of the wearer, around the base of the testicle-pouch.\nAdditionally hereto, in accordance with another suitable embodiment of the invention, one or more pre-tensioned elastic threads or the like may be arranged transversely of the shield, so as to tighten the shield around the penis and testicle-pouch of the wearer."}
{"text": "1. Field of the Invention\nThe present invention relates to a portable computer which can be easily carried and brought.\n2. Description of the Related Art\nA conventional portable computer is provided with a hinge portion at the front side spaced from the rear end of the body section by a certain distance (for example, U.S. Pat. No. 4,571,456). This is because a battery chamber and a modem enclosing portion are provided at the rear side of the hinge portion. Consequently a problem exists in that the dimensions of a cover section are limited and the dimensions of a display unit installed in the cover section can not be made large.\nAlso in a conventional portable computer, since a battery chamber and a modem enclosing portion are provided at the rear side of the hinge portion, these members limit the rotational angle of the cover section in an open state, relative to the body section, to about 135 degrees. The angle of 135 degrees does not produce a problem when the body section is used on a desk. However, if the computer has a liquid crystal display unit the body section is put and on a person's knee during the person's use outdoors or during use in a vehicle, the person's eyes are not orthogonal to the display unit and viewing of the display content becomes difficult.\nIn fact, from the viewpoint of human engineering, it is preferable that the open angle of the cover section is about 155 degrees, at least when the computer is to be placed for use on the user's knee. Consequently, a conventional portable computer is limited substantially to use on a desk rather than use on a knee."}
{"text": "Antennas that rely on the opening and closing of switches that are co-located with the antenna for tuning are well known in the prior art. An example of a MEMS tuned slot antenna used for frequency tuning is described in a co-pending U.S. Patent Application (See document number 1 below). The MEMS tuned slot antenna disclosed therein contains a slot that is shorted at one end and open at the other end, with a MEMS switch serving as the short across the open end, to determine the effective length of the slot. By closing different switches along the length of the slot, the frequency of the antenna can be tuned. At resonance, the slot measures one-half wavelength long from the closed end to the first closed MEMS switch. This antenna represents an improvement over previous tunable antenna designs because the current was forced through the switch due to the open end of the slot, thus eliminating any unwanted current paths through the ground plane. However, the effective size of this antenna is dependent on the wavelength, which can create problems when a compact antenna is needed. In general, to make any effective MEMS-tuned antenna, the MEMS switch should provide the only path for one part of the antenna current, because the finite inductance of the switch can be shorted by other nearby metal structures, particularly continuous ground planes.\nOther types of MEMS tuned antennas include patch designs, such as those described in document numbers 7 and 8 (identified below), as well as dipole, and various others. These designs are not preferred because patches, dipoles, and many other antennas are tuned by adding small metal regions that extend the length of the primary metal region. When tuning is performed with MEMS switches, this often causes interference from the DC bias lines. Therefore, it is necessary that the tuning be accomplished by shorting a metal object to a large ground plane, which can serve as both a RF and DC ground. In this way, the DC bias lines can be printed along this ground plane in such a way that they have very high or very low RF impedance, so that they cause minimal interference or coupling to the radiation. The slot antenna discussed above is an ideal candidate, but it suffers from a large size. It also requires that the ground plane be extended on all edges except one, which is left open for tuning.\nThus, the two important properties for a MEMS-tuned antenna are that the MEMS switch should be the only path for the particular portion of the antenna current that provides the tuning, and the switch should be able to be attached to a large ground plane to avoid interference or coupling from the DC bias. Another important property for many portable electronics or other compact devices is that the antenna should be small compared to the operating wavelength. One antenna that embodies these features is known as an F antenna. It typically consists of a metal wire or strip lying adjacent to the edge of a ground plane, with two connecting posts, one post acting as a feed for the metal strip, and the other acting as a short for impedance matching purposes. Reference 9 below discloses an F antenna by using a loop section for tuning instead of tuning the antenna itself. This design is not nearly as elegant or flexible, as the antenna does not provide a wide and arbitrary tuning range.\nThe disclosed antenna addresses the aforementioned needs by providing a simple, compact tunable antenna that is suitable for handheld or portable applications. The antenna can be tuned over a broad frequency range, and the size of the antenna is not solely dependent on the operating wavelength of the antenna such as is the case with typical prior art antennas.\n2. Description of Related Art                1. D. Sievenpiper, “RF MEMS-Tuned Slot Antenna and a Method of Making Same”, U.S. Patent Application Ser. No. 60/343,888 and U.S. patent application Ser. No. 10/192,986, which is related to 60/343,888. These applications describe a tunable slot antenna. The presently disclosed technology is different in that the presently disclosed technology allows an antenna to be much smaller than the operating wavelength which can be important for certain handheld and/or portable applications.        2. I. Korisch, “Planar Dual Frequency Band Antenna”, U.S. Pat. No. 5,926,139 describes a basic planar RF antenna and includes meander line type structures for setting the resonant frequency.        3. S. Moren, C. Rowell, “Trap Microstrip PIFA”, U.S. Pat. No. 6,380,895. This patent describes another type of planar RF antenna, and also includes meander line structures for setting the resonant frequency.        4. N. Johansson, “Antenna Device and Method for Portable Radio Equipment”, U.S. Pat. No. 6,016,125. This patent describes an antenna that is tunable or reconfigurable by adjusting the position of a whip portion, which contacts an impedance matching inductor. This could be used either to adjust the position of the antenna to improve the impedance match, or presumably to tune the resonant frequency of the antenna. However, this antenna requires physical control of the antenna position by a user, and the antenna is largely stationary.        5. Y. J. Chen, H. J. Li, R. B. Wu, “Multi-Resonance Horizontal U-Shaped Antenna”, U.S. Pat. No. 5,644,319. This patent describes a multi-resonant antenna, however the antenna is not tunable. Furthermore, the antenna requires a folded structure that increases the size of the antenna.        6. Hiroshi Okabe, Ken Take, “Tunable Slot Antenna with Capacitively Coupled Island Conductor for Precise Impedance Adjustment”, U.S. Pat. No. 6,034,655. This patent describes a slot antenna using a cavity structure. The cavity structure increases the size of the antenna significantly, and the use of a closed-end slot forbids the use of MEMS switches.        7. Robert Snyder, James Lilly, Andrew Humen, “Tunable Microstrip Patch Antenna and Control System Therefore”, U.S. Pat. No. 5,943,016 describes a method of using a patch antenna by using RF switches to connect or disconnect a series of tuning stubs. However, this antenna is extremely sensitive to the position of the bias circuits and does not have the ability to tune the polarization and the pattern.        8. Jeffrey Herd, Marat Davidovitz, Hans Steyskal, “Reconfigurable Microstrip Array Geometry which Utilizes Microelectromechanical System MEMS switches”, U.S. Pat. No. 6,198,438 describes an array of patch antennas that are connected by RF MEMS switches. This antenna can be selectively tuned by turning on or off various switches to connect the patches together. Larger or smaller clusters of patches will create antennas operating at lower or higher frequencies. However, this antenna requires a large number of switches and the antenna does not provide a way to eliminate the problem of interference between the DC feed lines and the RF part of the antenna.        9. Gerard Hayes, Robert Sadler, “Convertible Loop/Inverted F Antennas and Wireless Communicators Incorporating the Same”, U.S. Pat. No. 6,204,819 describes an F-type antenna. However, this antenna has significant drawbacks due to its complexity. The antenna requires each separate frequency of operation to be addressed by a different type of antenna (loop, F, etch). This requires a different set of design equations for different resonant frequencies and modes of operation. Furthermore, this antenna does not allow for angle diversity.        10. De Los Santos “Tunable Microwave Network Using Microelectromechanical. Switches” U.S. Pat. No. 5,808,527 describes a MEMS switch for tuning, but does not discuss integration of a switch into an antenna.        11. Lam, Tangonan, and Abrams, “Smart Antenna System Using Microelectromechanically Tunable Dipole Antennas and Photonic Bandgap Materials” U.S. Pat. No. 5,541,614 describes an antenna system using microelectromechanically tunable dipole antennas and photonic bandgap materials.        "}
{"text": "There exist many devices for extending and retracting a sheet of cloth or other material on a roller. These devices include window shades and rollers for screening fabric, such as mosquito netting. In general, these devices comprise a hollow roller having a winding spring connected to one side. The material is wound around the roller to take up the smallest possible space when the device is in the retracted position. In the extended position, the material is unwound from the roller, and the winding spring is under high tension. When the device is to be returned to the retracted position, tension on the winding spring is released causing the roller to rotate and the material to wind around the roller. Disadvantageously, however, the high load on the winding spring can cause the roller to rotate so fast during retraction that the device can cause injury to the operator, and can also cause loud unpleasant noises.\nTherefore, there is a need for a device for extending and retracting a sheet of material on a roller that is not associated with this disadvantage. Further, there is a need for a method of extending and retracting a sheet of cloth or other material on a roller that is not associated with this disadvantage."}
{"text": "Modern bathroom design favors a look with planer surfaces and configurations which minimize the visibility of unattractive space. Nowhere is this more noticeable than in the design of bathtubs and tub accessories. Traditional tubs had a shape matching their name, showing a curved exterior surface and exposed piping. Modern design provides more boxlike shape which hides piping and unattractive space beneath and around the bottom of the tub. A common way to provide this shape as part of a new tub system or often over the top of an older tub is to install a tub skirt. A tub skirt provides a flat, aesthetically pleasing front surface.\nA tub skirt, while hiding unattractive space, also prevents access to this space. As a remedy, tub skirts have been provided with access openings. These openings provide access for cleaning or tub repair and to access motors and plumbing on Jacuzzi type tubs. When such activities are not occurring, the openings are covered. Common accessories used to cover these openings are decorative panels. U.S. Pat. No. 5,940,906 to Halloran discloses a skirt frame with a detachable panel. The panel is held in place with Velcro tabs. U.S. Pat. No. 5,804,898 to Kapp et al. discloses a skirt frame with a detachable panel. Velcro is used for attachment. U.S. Pat. No. 5,208,924 to Smith et al. discloses a skirt frame and mechanically attachable panel. These panels are attached by screws with caps.\nSkirts are often provided without panels. In other cases, the panels provided with the skirts are not aesthetically pleasing and retrofitting is desired by the consumer. Common ways to attach aftermarket or retrofit panels in order to match new bathroom colors or the like is with velcro pads because they can be attached with an adhesive. However, velcro tends to lose holding force over time and cannot be cleaned easily. This lack of holding force is especially noticeable in Jacuzzi type tubs which include motors which vibrate the tub skirt slightly. Panels which are held in place with screws and caps are difficult to remove. Removal is a time consuming process. What is desired is a simple system which remains hidden from view, but can be used repeatedly without a loss in effectiveness."}
{"text": "Microfluidic liquid chromatography (LC) columns, or chips, packed with chromatographic particles require frits at an outlet and in most cases at an inlet of the column in order to retain a packed bed and ensure stable operation over time. In conventional LC columns with an inner diameter of 2.1 or 4.6 mm, for example, frits are typically made out of stainless steel or titanium particles sintered together, or screens, meshes and composites, that are mechanically press-fit into the ends of the column. In columns made out of fused silica capillaries, with inner diameters 75 to 150 microns, the frits are typically made with chemical means by dipping the end of the capillary into a polymeric solution that locks the ends of the column in place when it cures. Another method for capillary columns is to burn a section of the capillary to fuse the particles together.\nFrits can be made before or after packing. For example, it is possible to frit the outlet of the column and pack against it, then frit the inlet. Alternatively, a temporary external frit can be placed against the outlet of the column during packing, for example in the outlet capillary that carries the slurry liquid to waste, and removed after packing is complete, after which permanent frits are made on the capillary column at the outlet and possibly the inlet.\nIn planar microfluidic devices for LC applications, there can be several possible methods for fabricating frits. The geometry of such a device is shown schematically in FIG. 1. In this illustration, the device 100 is made by joining two layers. The top layer 101 contains two through-holes, or vias, 102. The lower layer 103 contains a groove 104. After alignment and joining of the two layers, a hermetically sealed channel is formed. Fittings are attached to the device 100 in order to make fluid connections to the vias 102. Particles dispersed in a slurry flow through the inlet via into the channel and out of the outlet via. A frit, either temporary or permanent, must be placed in the outlet via or after the outlet via in order to retain particles. If the frit is temporary and removed after packing is complete, a permanent outlet frit is made after packing. This frit must be capable of withstanding substantial force applied in subsequent operation without moving or rupturing. Frit motion would likely cause degradation of the column performance. Frit rupture would typically cause complete failure of the column. Typically, a permanent inlet frit is also created so that the packed bed is firmly locked into place.\nA first class of frits for microfluidic devices uses chemical solutions such as silicate. A drop is placed at the vias and allowed to cure, upon which the polymeric solution creates bridges between the particles and physically locks them in place. Another class of frits for microfluidic devices employs frits that are micro-machined along with other features of the device, such as channels and vias. In prior art techniques, the physical restriction that achieves retention of the particles inside the microfluidic device is an integral component of the device and is fabricated along with it. As such, the retaining device in those situations is fabricated prior to packing the particles into the microfluidic columns. Therefore, those retaining devices can only be fabricated at the column outlet.\nMechanical frits, typically made out of stainless steel or titanium particles compacted and sintered together, are a standard method for creating retaining structures in traditional LC columns. Typically, in a traditional column (see FIG. 2), made from a stainless steel tube 200, a first frit ring assembly 201, consisting of a frit 203 inside a frit ring 202, is placed at the outlet end of the tube 200 and maintained in place by an outlet end nut 204. Particles in a slurry are packed against this frit. A second frit assembly 205 is placed against the inlet of the tube 200 after packing and maintained in place using the inlet end nut 206."}
{"text": "Previously, an intake air control device, which opens and closes an intake air passage formed in a valve body by rotating a throttle valve received in the intake air passage, is widely used.\nIn a device, which is disclosed in the patent literature 1 as one such intake air control device, a valve shaft is fixed to the throttle valve at two axial locations by a pair of fixing members. With this fixing structure, rotatable shaft portions of the valve shaft, which respectively project from two opposite axial sides of the throttle valve, are rotatably supported by two separate bearings, respectively, which are held by the valve body. Among these two bearings, a first bearing, which rotatably supports the rotatable shaft portion located on one side, includes an inner race, which is fixed to the rotatable shaft portion located on the one side, and an outer race, which is fixed to the valve body. Furthermore, among the two bearings, a second bearing, which rotatably supports the rotatable shaft portion located on the other side, includes an inner race, which is axially slidably fitted to the rotatable shaft portion located on the other side, and an outer race, which is fixed to the valve body.\nIn the intake air control device disclosed in the patent literature 1, an aluminum-based material is used as a material of the valve body and the throttle valve, and an iron-based material is used as a material of the valve shaft. In this case, generally, a linear expansion coefficient of the throttle valve and a linear expansion coefficient of the valve body are larger than a linear expansion coefficient of the valve shaft in the axial direction. Therefore, the valve shaft, which is axially slidable at the second bearing side that is opposite from the first bearing having the inner race fixed to the valve shaft, will have relative displacement relative to the valve body due to a thermal expansion coefficient difference. Therefore, a recess for relieving the relative displacement around at least the shaft portion of the valve shaft located on the second bearing side is formed at an axial rim section of the throttle valve. In this way, even when the relative displacement of the valve shaft is generated, interference between the throttle valve and the valve body is limited, and thereby characteristics as of the intake air control device can be ensured.\nHowever, when the recess is formed at the throttle valve like in the disclosure of the patent literature 1, the intake air leaks through the recess even in the state where the throttle valve is fully closed. Therefore, at this time, the intake air flows through the intake air passage while the flow amount of this intake air in the intake air passage exceeds an inevitable flow amount of the intake air that inevitably flows through a full closing state clearance that is a clearance formed between the throttle valve in the full closing state and the valve body (hereinafter also simply referred to as a full closing state clearance). Thus, it is difficult to ensure the intended characteristics as of the intake air control device."}
{"text": "(1) Field of the Invention\nThe present invention relates to a method of manufacturing an optical filter adapted to selectively transmit a light of specific wavelength, and more particularly to a method of manufacturing an optical filter formed with dielectric multilayer coatings.\n(2) Related Art\nThere have been known a variety of optical filters for extracting a light of specific wavelength. As representative ones, such as those adopting a dielectric multilayer coating are utilized mostly. Generally, these optical filters adopting a dielectric multilayer coating are manufactured by alternately evaporating, onto a substrate glass, materials having different refractive indices such as TiO2 and SiO2, thereby forming a multiple interference coating (dielectric multilayer coating) by 50 to 100 layers. As the number of layers in the dielectric multilayer coating is increased, this type of optical filter is to have a relatively flat transmission characteristic and a sharply changing cutoff characteristic.\nFurther, as optical filters having a high resolving power, there have been known those elements such as called xe2x80x9cFabry-Perot interferometer (Fabry-Perot etalon, or simply etalon)xe2x80x9d. Generally, etalon is manufactured by applying coatings having high reflectance onto the side faces of a plane-parallel plate.\nMeanwhile, as a method of achieving large-capacitance of optical transmission system, attention has been directed to a wavelength-division multiplexing (WDM) optical transmission system in which two or more optical signals having wavelengths different from each other are multiplexed and transmitted in a single transmission path. In this WDM optical transmission system, since there are transmitted a plurality of WDM optical signals having different wavelengths arranged with very narrow wavelength intervals, there is required an optical filter having a high resolving power with a relatively flat transmission characteristic and a sharply changing cutoff characteristic, so as to select an optical signal of a specific wavelength from these WDM optical signals.\nHowever, although the aforementioned conventional optical filters adopting dielectric multilayer coatings are supposed to realize a flat transmission characteristic and a high cutoff characteristic, by increasing the number of the layers as mentioned above, the increase in the number of layers undergoes a predetermined limit such as due to restriction in actual manufacturing. Thus, it has been practically difficult to manufacture an optical filter adopting a dielectric multilayer coating, having both of a flat transmission characteristic and a high cutoff characteristic.\nFurther, the conventional etalons generally have a higher resolving power as compared to the optical filters adopting dielectric multilayer coatings. However, such etalons are problematic in that they have less flatness of transmission characteristic and have a deteriorated cutoff characteristic relative to adjacent bands. It has been also problematic that the yield of etalons having a high resolving power is low, since it has been difficult to form reflective coatings having uniform characteristics on the side faces of the plane-parallel plate.\nThe present invention has been carried out in view of the aforementioned points, and it is therefore an object of the present invention to provide a manufacturing method of readily realizing an optical filter having a flat transmission characteristic and a superior cutoff characteristic relative to adjacent bands.\nTo this end, with one embodiment of the present invention, there is provided a method of manufacturing an optical filter comprising two band-pass filters consisting of dielectric multilayer coatings, the optical surfaces of the filters being substantially parallelly arranged to each other, comprising: a dielectric multilayer coating forming step of forming a band-pass filter coating consisting of a dielectric multilayer coating, on one plane of a transparent substrate; a dividing step of dividing said substrate formed with said band-pass filter coating consisting of the dielectric multilayer coating, into a plurality of segments; and an adhering step of adhering those two segments which are selected from said plurality of segments of said substrate, by means of a transparent adhesive, so that respective surfaces of said band-pass filter coatings consisting of the dielectric multilayer coatings of said two segments are faced to each other and are rendered to be substantially parallel to each other with a predetermined distance.\nAccording to such a manufacturing method, the band-pass filter coating consisting of a dielectric multilayer coating is formed on the substrate in a single step, the substrate is divided, and the surfaces of the segments applied with the dielectric multilayer coating are adhered to manufacture an optical filter. This optical filter corresponds to what is obtained by forming band-pass filters instead of reflective coatings to be applied to a conventional etalon, thereby enabling provision of an optical filter having a flat transmission characteristic and a superior cutoff characteristic.\nIt is preferred that the adhering step in the abovementioned manufacturing method selects those adjacent two segments from the plurality of segments of the substrate divided by the dividing step. By selecting and adhering those adjacent segments of dielectric multilayer coatings formed in the same step, there are combined two band-pass filters having uniform characteristics.\nFurther, the abovementioned manufacturing method may further comprise an antireflection coating forming step of forming an antireflection coating on the other plane of the substrate, wherein the dividing step divides the substrate formed with the band-pass filter coating consisting of the dielectric multilayer coating and formed with the antireflection coating, into the plurality of segments. As a concrete example, the adhering step may adhere the selected two segments with a transparent thin plate interposed therebetween; and also it is preferable to set the distance between the band-pass filter coatings consisting of dielectric multilayer coatings of the selected two segments, to be equal to or less than 100 xcexcm.\nWith another embodiment of the present invention, there is provided a method of manufacturing an optical filter constituted of two band-pass filters consisting of dielectric multilayer coatings, the optical surfaces of the filters being substantially parallelly arranged to each other, comprising: a dielectric multilayer coating forming step of forming on a front surface of a transparent substrate a band-pass filter coating consisting of a dielectric multilayer coating rendered to have a grading characteristic concerning a transmission wavelength, and forming on a back surface of the transparent substrate a band-pass filter coating consisting of another dielectric multilayer coating rendered to have a grading characteristic in a direction perpendicular to that of the grading characteristic of the band-pass filter coating formed on the front surface; and a cutting step of cutting out a segment located at a specific area of the substrate formed with the band-pass filter coatings consisting of the dielectric multilayer coating at the front and back surfaces thereof, respectively.\nAccording to such a manufacturing method, two band-pass filter coatings consisting of two dielectric multilayer coatings having grading characteristics concerning transmission wavelengths are formed on the front and back surfaces of the substrate, respectively, such that the grading directions of the band-pass filter coatings are rendered to be perpendicular to each other. As a result, those segments of band-pass filters of uniform characteristics formed on the front and back surfaces, are distributed along a specific direction. Thus, there can be obtained an optical filter having an appropriate transmission characteristic, by cutting out a segment located at the specific area.\nPreferably, the dielectric multilayer coating forming step in the abovementioned manufacturing method forms the band-pass filter coatings consisting of dielectric multilayer coatings on the front and back surfaces of the substrate having a thickness of equal to or less than 100 xcexcm.\nFurther objects, features and advantages of the present invention will become more apparent from the following description of preferred embodiments when read in conjunction with the accompanying drawings."}
{"text": "EP-AI-No. 0059694 describes a drug administration device for releasing a drug into a flow of liquid during administration by a parenteral route to a patient, such device comprising a cell having an inlet and outlet for the liquid so as to permit the liquid to flow through the cell, the cell containing a solid drug-releasing unit positioned to contact the liquid flowing through the cell in use, the drug releasing unit being adapted to release the drug at a controlled rate into the liquid. Such device is gaining interest as a valuable new concept in intravenous infusion therapy i.a. The present inventors have found that by an improvement to such device substantial advantages can be gained, and further employment in extended therapeutical fields can be made possible.\nEP-AI-No. 0077604 and GB-A-No. 2107191 describe further devices for drug release to a liquid for parenteral administration."}
{"text": "1. Field of the Invention\nThe present invention relates to an integrated circuit, and more particularly to a high voltage integrated circuit.\n2. Description of the Prior Art\nA variety of power supplies and motor drivers utilize the bridge circuits to control a power source to the load. The bridge circuit normally has a high-side transistor connected to the power source and a low-side transistor connected to ground. A common node between the high-side transistor and low-side transistor is coupled to the load. As transistors are controlled to alternately conduct, the voltage of the common node swings between the voltage of the power source and the ground. Therefore the control of high-side transistor requires a charge pump circuit and/or a floating driver in order to fully turn on the high-side transistor. In recent development, many floating circuits are disclosed in U.S. Pat. No. 6,344,959 (Milazzo), U.S. Pat. No. 6,781,422 (Yang) and U.S. Pat. No. 6,836,173 (Yang).\nFIG. 1 shows a high-side transistor drive circuit, in which a circuit 10 is the floating driver. A capacitor 40 provides a supply voltage to the floating driver 10. A voltage VD charges the capacitor 40 through a diode 45 once the low-side transistor 30 is switched on. The ground reference of the capacitor 40 is pulled to the level of the voltage source VIN when the high-side transistor 20 is turned on. This happens, because turning on the high-side transistor 20 shifts the bridge circuit into a low impedance state."}
{"text": "In the packaging of bottled beer, there has recently been suggested in U.S. Pat. No. 3,815,808 to provide a packaging structure, especially for beer bottles, which consists of a substantially rigid tray which may be reusable several times and typically constructed of plastic material and a disposable sleeve which closes the open top of the tray. The disposable sleeve includes downwardly depending panels integrally formed with a top panel coextensive with the open top of the tray, the downwardly depending panels extending substantially the height of the tray walls and gripping the tray walls thereby inhibiting the removal of the sleeve from the tray. The top panel of the sleeve has an opening means whereby the consumer may gain access to the beer bottles in the tray, the sleeve, however, remaining in association with the tray after opening of the opening means.\nThe packaging of beer bottles for sale is a high speed generally automated operation requiring precise and reliable machinery. Currently, there is not available packaging machinery suitable for high speed operation in the formation of packaging structures in accordance with the above-mentioned U.S. Pat. No. 3,815,808."}
{"text": "1. Field of the Invention\nThe present invention relates to haptic seats of the type used in vehicles, such as automobiles\n2. Background Art\nA haptic vehicle seat is configured to impart sensations to an occupant thereof. The sensation may include vibrations, pulses, shakes, or some other inducement that involves the occupants sense of touch."}
{"text": "1. Field of the Invention:\nThis invention relates to matrix display devices, that is to display devices comprising a large plurality of elements arranged adjacently and each individually controllable to display different aspects to an observer.\n2. Description of the Prior Art\nThere are known in the art from United Kingdom patent specifications Nos. 355,440 and 934,001 electrically controlled dot matrix display devices in which each dot element includes a laminar member, having differently reflective opposed faces, arranged to be rotatable about an axis parallel to the plane of the array so as to build up characters out of selections of reflective dots.\nFor night viewing, it is necessary to illuminate the prior art displays from the front, using relatively remote light sources."}
{"text": "The present invention relates to a salt and associated hydrates of racemic 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid, processes for their preparation, pharmaceutical compositions comprising them, and their use in antibacterial therapy.\nEP 688772 (corresponding to Korean Patent Laid open Publication No 96-874) discloses novel quinoline(naphthyridine)carboxylic acid derivatives, including anhydrous 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid of formula I, having antibacterial activity. \nAccording to the invention there is provided 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate.\n7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate (hereinafter referred to as xe2x80x98the methanesulfonatexe2x80x99) may be obtained as an anhydrate or a hydrate (i.e., methanesulfonate.nH2O).\nHydrates of the methanesulfonate wherein n is in the range of from 1 to 4 are preferred. Particular hydrates of the methanesulfonate, which may be mentioned, are those in which n is 1, 1.5, 2, 2.5, 3, 3.5, and 4. Particularly preferred compounds are those in which n is 1.5 or 3, with n=1.5 being most preferred.\nThe moisture content of the methanesulfonate hydrates varies with the hydration number (n) of the hydrated molecule. The methanesulfonate has a molecular weight of 485.5. Thus the calculated moisture content of hydrates where n is 1, 1.5, 2, 2.5, 3, 3.5, and 4 is 3.6%, 5.0%, 6.9%, 8.5%, 10.0%, 11.5%, and 12.9%, respectively. However, the actual moisture content of the methanesulfonate hydrates may differ from the calculated value depending on various factors including recrystallization conditions and drying conditions. The observed moisture content for the methanesulfonate hydrates where n is 1, 1.5, 2, 2.5, 3, 3.5, and 4 is shown in Table 1:\nIt is possible to mix methanesulfonate hydrates having different moisture contents together to give materials having intermediate moisture contents.\nPreferred methanesulfonate hydrates have a moisture content of from 4 to 6% or from 9 to 11%, especially a moisture content of from 4 to 6%.\nThe methanesulfonate has been observed to exist as a stable hydrate over a range of hydration numbers (n). Stability of the hydrate refers to its resistance to loss or gain of water molecules contained in the compound. The methanesulfonate hydrates maintain a constant moisture content over an extended relative humidity range. The n=3 hydrate has a constant moisture content at a relative humidity of from at least 23 to 75%, and the n=1.5 hydrate has a constant moisture content at a relative humidity of from 23 to 64% (see FIGS. 3 and 4). In contrast, moisture absorption by the anhydrate varies greatly with relative humidity.\nBoth the methanesulfonate anhydrate and n=3 hydrate undergo transition to the n=1.5 hydrate in aqueous suspension indicating that the latter is thermodynamically more stable. The n=1.5 hydrate is a sesquihydrate at 11 to 64% of relative humidity. Above 75% relative humidity, it takes up water over 10% and its X-ray diffraction pattern changes. The hydrate (another form of n3 having different physicochemical properties from the n=3 hydrate of Example 2) obtained from n1.5 hydrate at 93% relative humidity is not stable at lower relative humidity, and it converts back to n=1.5 hydrate at a relative humidity below 75%.\nSince the moisture content of the anhydrate changes readily depending on the environment (e.g., relative humidity, formulation additives, etc.), it may require careful handling during storage or formulation, with operations, such as quantifying procedures, being performed in a dry room. The hydrates do not change in moisture content easily and hence products, which are stable under prolonged storage and formulation may be obtained. The hydrate can be tableted without the addition of a binder since the water contained in the compound itself acts as a binder, whereas it may not be possible to tablet the anhydrate at a similar pressure.\nThe present invention also provides a process for the preparation of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate and hydrates thereof which comprises reacting 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid with methanesulfonic acid and crystallizing the resulting methanesulfonate from solution, and where desired or necessary adjusting the hydration of the compound.\nThe methanesulfonate and its hydrates may be prepared by the addition of methanesulfonic acid to the free base which may be prepared as described in EP 688772. Preferably, 0.95 to 1.5 molar equivalents of methanesulfonic acid are added to the free base, or 1 molar equivalent of methanesulfonic acid dissolved in a suitable solvent is added to the free base. Suitable solvents for the preparation of the methanesulfonate and its hydrates include any solvent in which the methanesulfonate is substantially insoluble, and the suitable solvents include C1-C4 haloalkanes, C1-C8 alcohols and water, or mixtures thereof. Dichloromethane, chloroform, 1,2-dichloroethane, methanol, ethanol, propanol and water, or mixtures thereof, are preferred solvents. If necessary, the free base may be heated in the solvent to facilitate solution before methanesulfonic acid is added, and alternatively the methanesulfonic acid may be added to a suspension, or partial suspension, of the free base in the solvent. Following addition of the methanesulfonic acid, the reaction mixture is preferably allowed to stand or is stirred for 1 to 24 hours at a temperature of from about xe2x88x9210 to 40xc2x0 C. The resulting methanesulfonate is obtained as a solid, which can be isolated by filtration or by removal of the solvent under reduced pressure.\nDifferent hydrates may be obtained by altering the recrystallization conditions used in the preparation of the methanesulfonate, and such conditions may be ascertained by conventional methods known to those skilled in the art.\nThe present invention also provides a process for the preparation of a hydrate of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate comprising exposing the methanesulfonate anhydrate or a solvate thereof to a high relative humidity.\nThe methanesulfonate anhydrate or solvate thereof is preferably exposed to a relative humidity of at least 75%.\nThe methanesulfonate anhydrate or solvate thereof may be exposed to high relative humidity by passing humidified nitrogen gas through the methanesulfonate anhydrate or solvate thereof or by standing the methanesulfonate anhydrate or solvates thereof under a high relative humidity.\nThe humidified nitrogen gas used in this process, for example nitrogen gas having a humidity of at least 75%, may be made by conventional methods. In this process it is desirable to maintain the temperature in the range above which moisture condensation could occur. Also, particularly in large scale production, it is preferable to stir the sample thoroughly while the humidified nitrogen gas is passed through. When the hydrate is prepared by standing the methanesulfonate anhydrate or solvate thereof under a high relative humidity, for example a relative humidity of at least 75%, it is preferable to spread the sample as thinly as possible in order to raise the conversion efficiency.\nThe solvates of methanesulfonate anhydrate which may be used in the process according to this aspect of the present invention include solvates with one or more organic solvents. Preferred solvents include C1-C4 haloalkanes and C1-C8 alcohols, for example those selected from the group consisting of ethanol, dichloromethane, isopropanol and 2-methyl-2-propanol.\nSolvates of the methanesulfonate anhydrate are novel. Thus according to a further aspect of the invention there is provided a solvate of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate with one or more organic solvents.\nThe solvates of the methanesulfonate are prepared by recrystallization and controlled by the condition of recrystallizing system.\nThe methanesulfonate and its hydrates exhibit the same potent antibacterial activity as the corresponding free base disclosed in EP 688772. The methanesulfonate and its hydrates also exhibit desirable physicochemical properties including improved solubility and constant moisture content regardless of the ambient relative humidity when compared to the free base and other salts thereof. The methanesulfonate and its hydrates thus exhibit greater ease of handling, quality control and formulation than the free base and other salts thereof.\nAs mentioned above the methanesulfonate and its hydrates exhibit antibacterial activity. The methanesulfonate and its hydrates may be formulated for administration in any convenient way for use in human or veterinary medicine, according to techniques and procedures per se known in the art with reference to other antibiotics, and the invention therefore includes within its scope a pharmaceutical composition comprising 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate or a hydrate thereof together with a pharmaceutically acceptable carrier or excipient.\nCompositions comprising the methanesulfonate or hydrate thereof as active ingredient may be formulated for administration by any suitable route, such as oral, parenteral, or topical application. The compositions may be in the form of tablets, capsules, powders, granules, lozenges, creams, or liquid preparations, such as oral or sterile parenteral solutions or suspensions. Tablets and capsules for oral administration may be in unit dose presentation form and may contain conventional excipients, such as binding agents, for example, hydroxypropyl methyl cellulose, hydroxypropyl cellulose, syrup acacia, gelatin, sorbitol, tragacanth, or polyvinylpyrrolidone; fillers, for example, microcrystalline cellulose, lactose, sugar, maize-starch, calcium phosphate, sorbitol, or glycine; tableting lubricants, for example, magnesium stearate, talc, polyethylene glycol, or silica; disintegrants, for example, sodium starch glycolate, cross-linked polyvinylpyrrolidone, or potato starch; or acceptable wetting agents such as sodium lauryl sulfate. The tablets may be coated according to methods well known in normal pharmaceutical practice. Oral liquid preparations may be in the form of, for example, aqueous or oily suspensions, solutions, emulsions, syrups or elixirs, or may be presented as a dry product for reconstitution with water or other suitable vehicle before use. Such liquid preparations may contain conventional additives such as suspending agents, for example, sorbitol, methyl cellulose, glucose syrup, gelatin, hydroxyethyl cellulose, carboxymethyl cellulose, aluminum stearate gel or hydrogenated edible fats; emulsifying agents, for example, lecithin, sorbitan monooleate, or acacia; non-aqueous vehicles (which may include edible oils), for example, almond oil, oily esters, glycerine, propylene glycol, or ethyl alcohol; preservatives, for example, methyl or propyl p-hydroxybenzoate or sorbic acid; and, if desired, conventional flavoring or coloring agents. Suppositories will contain conventional suppository base, e.g., cocoa-butter or other glyceride.\nFor parenteral administration, fluid unit dosage forms are prepared utilizing the compound and a sterile vehicle, water being preferred. The methanesulfonate or hydrate thereof, can be either suspended or dissolved in the vehicle, depending on the vehicle and concentration used. In preparing solutions the methanesulfonate or hydrate thereof can be dissolved in water for injection and filter sterilized before filling into a suitable vial or ampoule and sealing. Advantageously, agents such as local anaesthetic, preservative and buffering agents can be dissolved in the vehicle. To enhance the stability, the composition can be lyophilized and the dry lyophilized powder sealed in a vial, and an accompanying vial of water for injection may be supplied to reconstitute the powder prior to use. Parenteral suspensions are prepared in substantially the same manner except that the methanesulfonate or hydrate thereof is suspended in the vehicle instead of being dissolved and sterilization cannot be accomplished by filtration. The methanesulfonate or hydrate thereof can be sterilized by exposure to ethylene oxide before suspending it in the sterile vehicle. Advantageously, a surfactant or wetting agent is included in the composition to facilitate uniform distribution of the methanesulfonate or hydrate thereof.\nThe methanesulfonate or hydrate thereof may also be formulated as an intramammary composition for veterinary use.\nThe composition may contain from 0.1% to 100% by weight, preferably from 10 to 99.5% by weight, more preferably from 50 to 99.5% by weight of the active ingredient measured as the free base, depending on the method of administration. Where the compositions comprise dosage units, each unit will preferably contain from 50-1500 mg of the active ingredient measured as the free base. The dosage as employed for adult human treatment will preferably range from 100 mg to 12 g per day for an average adult patient (body weight 70 kg), for instance 1500 mg per day, depending on the route and frequency of administration. Such dosages correspond to approximately 1.5 to 170 mg/kg per day. Suitably the dosage is from 1 to 6 g per day.\nThe daily dosage is suitably given by administering the active ingredient once or several times in a 24-hour period (e.g., the active ingredient up to 400 mg may be administered once a day). In practice, the dosage and frequency of administration which will be most suitable for an individual patient will vary with the age, weight, and response of the patients, and there will be occasions when the physician will choose a higher or lower dosage and a different frequency of administration. Such dosage regimens are within the scope of this invention.\nThe present invention also includes a method of treating bacterial infections in humans and animals comprising administering a therapeutically effective amount of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate or a hydrate thereof.\nIn a further aspect, the present invention also provides the use of 7-(3-aminomethyl-4-methoxyiminopyrrolidin-1-yl)-1-cyclopropyl-6-fluoro-4-oxo-1,4-dihydro-1,8-naphthyridine-3-carboxylic acid methanesulfonate or a hydrate thereof for the manufacture of a medicament for treating bacterial infection.\nThe methanesulfonate and its hydrates are active against a broad range of Gram-positive and Gram-negative bacteria, and may be used to treat a wide range of bacterial infections including those in immunocompromised patients.\nAmongst many other uses, the methanesulfonate and its hydrates are of value in the treatment of skin, soft tissue, respiratory tract and urinary tract infections, and sexually transmitted diseases in humans. The methanesulfonate and its hydrates may also be used in the treatment of bacterial infections in animals, such as mastitis in cattle."}
{"text": "The compound cyclosporine (cyclosporin A) has found wide use since its introduction in the fields of organ transplantation and immunomodulation, and has brought about a significant increase in the success rate for transplantation procedures. Unsatisfactory side-effects associated with cyclosporine, however, such as nephrotoxicity, have led to a continued search for immunosuppressant compounds having improved efficacy and safety.\nRecently, several classes of macrocyclic compounds having potent immunomodulatory activity have been discovered. Okuhara et al., in European Patent Application No. 184162, published Jun. 11, 1986, disclose a number of macrocyclic compounds isolated from the genus Streptomyces. Immunosuppressant FK-506, isolated from a strain of S. tsukubaensis, is a 23-membered macrocyclic lactone represented by formula 1a, below. Other related natural products, such as FR-900520 (1b) and FR-900523 (1c), which differ from FK-506 in their alkyl substituent at C-21, have been isolated from S. hygroscopicus yakushimnaensis. Yet another analog, FR-900525, produced by S. tsukubaensis, differs from FK-506 in the replacement of a pipecolic acid moiety with a proline group.\nFR-900520, also known as ascomycin, has been previously disclosed by Arai et al. in U.S. Pat. No. 3,244,592, issued Apr. 5, 1966, where the compound is described as an antifungal agent. Monaghan, R. L., et al., on the other hand, describe the use of ascomycin as an immunosuppressant in European Patent Application No. 323865, published Jul. 12, 1989.\nAlthough the immunosuppressive activity of FK-506 has been clinically confirmed, toxicity in mammals has limited its utility. The activity of FK-506 has, however, prompted efforts to discover novel analogs of FK-type compounds which possess superior properties. These efforts include the isolation of new fermentation products, the microbial transformation of existing chemical entities, the chemical modification of these macrocycles, and the synthesis of hybrid species derived from smaller synthetic fragments. ##STR3##\nFermentation products of FK-type compounds include C-21-epi derivatives of FK-506; a 31-demethylated derivative of FK-506; 31-oxo-FK-506; and compounds derived from FK506, FR-900523, and FR-900525 which are characterized by the introduction of hydroxy-protecting groups, formation of a double bond by elimination of water between carbons 23 and 24, oxidation of the hydroxy group at carbon 24 to the ketone, and reduction of the allyl side-chain at carbon 21 via hydrogenation. Other published metabolites include those derived from FK-506 and FR-900520 where the lactone ring is contracted to give a macrocyclic ring containing two fewer carbons.\nSeveral microbial transformations of FK-type compounds at carbon 13 have been published, such as the microbial demethylation of FR-900520 to form the bisdemethylated 13,31-dihydroxy ting-arranged derivative of FR-900520; the microbial monodemethylation of FK-506 and FR-900520, respectively; and the microbial demethylation of FR-900520 at C-31, as well as a number of other macrocyclic microbial transformation products.\nNumerous chemical modifications of the FK-type compounds have been attempted. These include the preparation of small synthetic fragments of FK-type derivatives; a thermal rearrangement of a variety of derivatives of FK-506 which expands the macrocyclic ring by two carbons; and modifications which include methyl ether formation at C-32 and/or C-24, oxidation of C-32 alcohol to the ketone, and epoxide formation at C-9.\nAlthough some of these modified compounds exhibit immunosuppressant activity, the need remains for macrocyclic compounds having immunosuppressive activity which do not have the serious side effects frequently associated with immunosuppressant therapy. Accordingly, one object of the invention is to provide novel semisynthetic macrolides which possess the desired immunomodulatory activity but which may be found to minimize untoward side effects.\nAnother object of the present invention is to provide synthetic processes for the preparation of such compounds from starting materials obtained by fermentation.\nA further object of the invention is to provide pharmaceutical compositions containing, as an active ingredient, one of the above compounds. Yet another object of the invention is to provide a method of treating a variety of disease states, including post-transplant tissue rejection and autoimmune disfunction."}
{"text": "The present invention relates to a continuous polymerization method of laurolactam (dodecanoic lactam) and the apparatus therefor. More specifically, the present invention relates to a continuous polymerization method of laurolactam (dodecanoic lactam), comprising continuously polymerizing laurolactam as a raw material monomer by using water as a catalyst, thereby producing a desired high quality Nylon 12 with a wide range of molecular weight, and the apparatus therefor.\nNylon 12 has conventionally been produced by heating and polymerizing laurolactam (dodecanoic lactam) or 12-aminododecanoic acid. By the method using 12-aminododecanoic acid as a raw material monomer, 12-aminododecanoic acid is continuously heated and melted under atmospheric pressure for polycondensation. Alternatively, the ring opening reaction of laurolactam progresses very slowly, which rate is 1/10 or less that of the rate of hydrolysis and polymerization of relatively easily ring opening caprolactam under atmospheric pressure as a raw material of Nylon 6. Therefore, the process using laurolactam as a raw material monomer generally comprises effecting a prepolymerization process wherein laurolactam is firstly heated along with water or an acid catalyst or the like under pressure, thereby promoting the ring opening reaction following a high-temperature/high-pressure system, rendering the pressure to atmospheric pressure, and polycondensing the resulting product under atmospheric pressure or reduced pressure to prepare a high molecular weight Nylon 12.\nIf it is intended to produce Nylon 12 by continuous polymerization using laurolactam as a raw material monomer, a large-scale, complex apparatus for high temperatures and high pressures is therefore inevitable due to the prepolymerization process, causing problems such as high equipment cost. Thus, a batch polymerization system (discontinuous two-step polymerization system) has conventionally been employed primarily for industrial polymerization of laurolactam.\nHowever, such a batch system has the following drawbacks; the operational procedure is complex with a low productivity; when drawing out a generated polymer after the termination of the post polymerization process, the adhesion of the generated polymer onto the polymerization vessel and the retention thereof in the vessel is unavoidable; if the prepolymerization and post polymerization are repeated in the same container, the residual polymer is deteriorated through thermal hysteresis, thereby reducing the quality of the product polymer. In order to prevent these drawbacks, the vessel should necessarily be washed periodically, resulting in a far lower productivity.\nIn producing Nylon 12 having a high polymerization degree from laurolactam as the raw material monomer, the most serious problem in either one of the continuous polymerization or the batch polymerization is a laurolactam polymerization-conversion rate at the prepolymerization process under pressure. For the ring opening polymerization of laurolactam, the conversion rate obtained by the prepolymerization process influences the final polymerization product. While the prepolymerization process is conducted under pressure, the procedures in the prepolymerization process thereafter are conducted under atmospheric pressure or reduced pressure, and because it cannot be expected to increase the conversion rate at the prepolymerization process by a much greater degree though it is sometimes possible that the rate slightly increase. If the prepolymerization process is insufficient, disadvantages may develop such as the occurrence of contamination of a nozzle or a mold during the molding process, so it is important for laurolactam polymerization that the polymerization-conversion rate (ring opening reaction) at the prepolymerization process under pressure should be facilitated to a nearly equilibrium value (about 99.7%).\nIn order to promote the conversion rate at the prepolymerization process up to at the equilibrium value in such case, water conventionally employed as a single catalyst is generally added to laurolactam at 1 to 10% by weight under a pressure of about 10 to 50 kg/cm.sup.2 G and a higher temperature above 300.degree. C. But at a higher temperature above 300.degree. C., thermal deterioration, namely, partial gelation or discoloring or the like, is likely to occur due to the occurrence of undesirable secondary reactions, so that a high-quality polymer cannot be obtained. At a temperature below 300.degree. C., alternatively, a prepolymerization time period of at least 10 hours or more is required, resulting also in the deterioration of the quality of a polymer due to the factor of thermal deterioration via a prolonged polymerization at a high temperature and a high pressure.\nIf the water addition rate is less than 1% by weight, the ring opening polymerization of laurolactam is extremely difficult to occur, whereby the polymerization addition rate is lowered. On the other hand, if above 10% by weight, the dehydration from a reaction product occurs involving bubble, so that the plug flowability of a reaction solution inside a reaction column is lowered, resulting in the variation of the conversion rate, unfavorably.\nFurthermore, when intending a high conversion rate by a prepolymerization at a temperature below 300.degree. C. and for 10 hours or less, acid or alkali catalysts are required besides water. When these catalysts are used, the molecular weight is hard to adjust at post polymerization, and hence, it is difficult to obtain a Nylon 12 with a desired molecular weight. Also, the viscosity change and the like at the molding process (at remelting) are remarked as problems, so that the industrial practice thereof is difficult.\nAs has been described above, the method comprising continuously polymerizing laurolactam employing water as a catalyst thereby industrially producing a high-quality Nylon 12 with a predetermined molecular weight has a variety of problems, so further improvement thereof has been desired.\nIt is the object of the present invention to overcome the problems in polymerizing laurolactam as a raw material monomer, and to provide a method for continuous polymerization of laurolactam, comprising industrially producing a high-quality Nylon 12 with a desirable molecular weight at a lower polymerization temperature and for an extremely short period of time, and the apparatus therefor."}
{"text": "It is known that hormones like the glucocorticoid and thyroid hormones enter cells by facilitated diffusion. It is also known that hormones then bind to specific receptor proteins, thereby creating a hormone/receptor complex. The binding of hormone to the receptor initiates an alosteric alteration of the receptor protein. As a result of this alteration, it is believed that the hormone/receptor complex is capable of binding with high affinity to certain specific sites on the chromatin DNA. Such sites, which are referred to as hormone response elements or HRE's, modulate expression of nearby target gene promoters.\nA major obstacle to further understanding of the specifics of gene regulation by exogenous inducers such as hormones has been the lack of availability of receptor proteins in sufficient quantity and purity to allow such proteins to be adequately analyzed and characterized. This same lack of availability has thwarted the use of receptors in diagnostic assays to determine the presence of exogenous inducers (e.g., the hormones) in various body fluids and tissues, as well as their use as \"prototypes\" for engineering chimeric receptor protein analogs.\nIn an effort to overcome this lack of availability of receptor proteins, scientific investigators are working to discover the genes that encode such proteins. To date several such genes have been disclosed and characterized. The cloned genes include those encoding the following receptors: glucocorticoid, mineralocorticoid, progesterone, estrogen, the two steroid-related receptors (known in the art as ERR1 and ERR2), vitamin D.sub.3, thyroid, v-erb-A, E75 (Drosophilia and two retinoid receptor proteins, retinoic acid receptor alpha (RAR.alpha.) and retinoic acid receptor beta (RAR.beta.). (See Giguere, et al., (1987) regarding RAR.alpha., and Petkovich, et al., (1987) and Brand, et al., (1988) regarding RAR.beta..)\nThis disclosure describes the isolation and characterization of a cDNA encoding a third functional retinoid receptor protein that is referred to herein as the gamma retinoic acid receptor (RAR.gamma.). Like RARs alpha and beta, the new gamma retinoic acid receptor has homology with the DNA-binding and ligand-binding domains of the steroid and thyroid hormone receptors.\nThe retinoic acid receptor genes belong to the superfamily of genes known as the steroid hormone receptor family. All genes in this family can be divided into discrete regions or domains that are sometimes referred to as regions A/B, C, D, E, and F. See FIG. 2; also see Robertson, (1987) and Evans, (1988). The C region encodes the DNA-binding domain, the E region encodes the ligand-binding domain and the F region encodes the carboxy-terminus domain. The D region is believed to function as a \"hinge\". The function of the A/B (or N-terminus) region is not entirely clear; it may be involved with enhancement and repression of receptor transcription activity. See for example, Hollenberg, et al., (1988) and Oro, et al., (1988).\nThe present specification also discloses chimeric receptors made by \"swapping\" functional domains between the new gamma retinoic acid receptor and the glucocorticoid, the mineralocorticoid, the progesterone, the estrogen, the estrogen-related (ERR1 and ERR2), the vitamin D.sub.3 receptor, the thyroid receptors, the v-erb-A receptor, the E75 (Drosophilia) receptor and the alpha and beta retinoic acid receptors. These chimeric receptors have hybrid functional characteristics based on the \"origin\" of the \"parental\" DNA-binding and ligand-binding domains incorporated within the chimeras. For example, if the DNA-binding domain in the chimeric receptor is the gamma retinoic acid receptor DNA-binding domain (i.e., is obtained from wild-type gamma retinoic acid receptor or is a mutant that contains the functional elements of the gamma retinoic acid DNA-binding domain), then the chimera will have DNA-binding properties characteristic of the gamma retinoic acid receptor. The same is true of the ligand-binding domain."}
{"text": "Often there is an online set of credentials that identify a user and match them to a piece or set of data that a third party would like to have that user matched to. The present invention is a method to use the success or failure of an online log-in using a set of credentials to create or verify a user on another system and/or match the user on the to a piece or set of data on a third system.\nThe invention has many uses across many industries, but will be discussed with regard to a specific example. There is a specific application of the invented method that solves an acute need in the world of online banking. Financial institutions (including banks, credit unions and other institutions that hold or process money for customers) have recognized the desire of customers or members to access their accounts, view financial data, and perform transactions online from a remote location such as from a home or office rather than in-person in a physical bank, or credit union, stock brokerage, etc. Customers, members and account holders (hereinafter “customers”) often find online banking more convenient and efficient than having to visit a branch in person and the number that find online banking more convenient is increasing.\nTo facilitate online banking, financial institutions (hereinafter “FI”) establish their own or contract with an online banking provider (OLBP) which allows the FI's customers to access their accounts and perform transactions through a remote computer system rather than in-person at the FI's physical location. Referring to FIG. 1, an example of such a system is depicted.\nIn FIG. 1, the FI 101 typically creates and manages all accounts as well as processes all transactions through or using the FI's core computer system (or core banking system) 102. That is where confidential data such as customer or member account information is stored, such as on a data store 102a. An OLBP 103 is established to provide an online banking interface with the FI's customers. The OLBP utilizes custom APIs to establish its access to information on the FI's core computer system. The OLBP 103 will typically store or connect to certain information, like usernames and passwords that allow a user to access online banking; on a database it controls or has access to.\nIn the prior art in a traditional OLBP arrangement, an FI's customers have accessed their financial accounts held at the individual's bank or credit union through the institution's OLBP 103 via a remote desktop computer 109. That computer may be located at a home, office or elsewhere. Such interaction was typically conducted through the internet or other network. The OLBP utilizes a web server to maintain such service.\nWith the growing popularity and ubiquity of mobile computing devices (hereinafter “MCD” 106), such as smart phones (Apple iPhone® and Android® phone), tablet computers (Apple iPad®, etc.), and other mobile devices, FI customers now wish to access their accounts and conduct their financial transactions wirelessly from those MCDs. However due to size constraints, device limitations, the strong preference for native apps, etc., the standard OLB experience does not normally work well on these MDCs.\nMobile banking providers (hereinafter “Mobile Banking Provider” or “MBP” 108 with its own data store 108a) arose to bridge this gap in user experience and develop sites that can be viewed better on the smaller screen of MCDs, as well as create native apps that work on the MCDs. In order for a user to authenticate to an MCD and gain access to their banking information, the MCD would normally need to have a connection to the OLBP that would allow the authentication to occur via an API or similar connection that involves some work or authorization from the OLBP for the MBP to connect in that way. Not all OLBPs are willing to allow this type of connection to MBPs and many charge burdensome or prohibitive fees to do so. The OLBP can also take a very long time to do the work/authorization on their side to allow the MBP to connect this way. Some relationships that MBPs desire to form with OLBPs may never come about, which may result in the financial institution not being able to provide certain mobile account functionality to its customers. Consequently an MBP investing in an application to support mobile banking on a particular MCD cannot be certain when, if ever, the application in question will be able to integrate with a particular FI. And with approximately 13,000 banks and credit unions in the United States alone, the task of creating necessary integrations with each financial institution's OLBP is enormous.\nIt is important to note that MBP are not alone in this, any provider wishing to provide services to a FI that wants users' accounts behind the OLBP wall would usually need to pair its users to OLBP user data behind that wall.\nThe limitations of this situation in the prior art should be obvious. Not only is it time-consuming and expensive to create the necessary relationships with the OLBPs and write all of the APIs, but those OLBPs must be maintained at great expense. The MBP also has no guarantee that the OLBP or similar partner will ever grant the authorization or do the work needed for the MBP to connect in this manner.\nFurther, as FIs generally focus on the banking business and not on the software development business or even on new technologies such as mobile devices and mobile banking, it is typically left to third party application developers and MBPs to develop relevant new software applications.\nMany FIs are under arduous and lengthy contracts with OLBPs that are not motivated to integrate best of breed new technologies, especially those the OLBP feels are a threat. In these situations the OLBPs often act like a “troll” that will either inflict punishing integration fees, costing the FI more than the MBP service in the first place to try to keep a competitor out. Often, the OLBP will unfortunately do this even if it is against the FI's wishes and best interest and even though the FI is supposed to be the customer the OLBP is meant to serve. This is a very abusive practice by the OLBP, the entrenched provider. Making matters even worse for the FI or the customer that is meant to be served, some OLBPs will simply not allow anyone to cross that integration bridge even if the FI or MBP is willing to pay the abusive fee.\nAt best the process of a MBP having to integrate with an OLBP in order to serve an FI is slow, inefficient and costly. At worst, it is either impossibly expensive or simply blocked, to the detriment of the FI.\nThe prior art reveals a clear need for a way to provide efficient connectivity between MCDs and FIs other than either working one at a time through OLBPs or establishing a connection direction with each of the many FIs. The invention seeks to provide a more efficient method or technique for MBP's to authenticate account information in order to correlate with financial institutions banking platforms and support mobile banking, PFM and other applications of that nature for all or substantially all financial institutions."}
{"text": "The present invention relates to a driving method for a liquid crystal device for use in flat-panel displays, projection displays, printers, etc.\nAs a type of a liquid crystal (liquid crystal device) for displaying various data (information) by using a liquid crystal, there have been known those using a nematic liquid crystal or a chiral smectic liquid crystal. A liquid crystal device using a chiral smectic liquid crystal has an advantage of, e.g., higher response speed than that using a nematic liquid crystal, thus being expected to be widely utilized.\nMore specifically, a twisted nematic (TN) liquid crystal has widely been used conventionally as a material for a liquid crystal device as described by M. Schadt and W. Helfrich, xe2x80x9cApplied Physics Lettersxe2x80x9d, Vol.18, No.4 (Feb. 15, 1971), pp. 127-128. The TN liquid crystal is used in an active matrix-type liquid crystal device (panel) in combination with switching elements such as thin film transistors (TFTs). The active matrix-type liquid crystal device is free from a problem of cross-talk and is produced with high productivity with respect to that having a size (diagonal length) of 10-17 in. with a progress of production technique.\nHowever, the above-mentioned liquid crystal device using the TN liquid crystal has been accompanied with problems such as a slower response speed and a narrower viewing angle.\nIn order to solve the problems, various alignment modes including an optically compensated bend or birefringence (OCB) mode for improving a response speed, and In-Plain Switching mode and Vertical Alignment mode for improving a viewing angle have been proposed but are not said to be satisfactory for improvements in response seed and/or viewing angle.\nIn order to solve the problems of the conventional TN liquid crystal devices, a liquid crystal device using a chiral smectic liquid crystal exhibiting bistability has been proposed by Clark and Lagerwall (Japanese Laid-Open Application (JP-A) 56-107216, U.S. Pat. No 4367924). As the liquid crystal exhibiting bistability, a ferroelectric liquid crystal having chiral smectic C phase is generally used. Such a ferroelectric liquid crystal provides a very quick response speed because it causes inversion switching of liquid crystal molecules based on their spontaneous polarizations. In addition, the ferroelectric liquid crystal assumes bistable state showing a memory characteristic and further has an excellent viewing angle characteristic, thus being considered to be suitable for a display device or light-valve of high speed, high definition and larger area.\nIn recent years, an anti-ferroelectric liquid crystal exhibiting tristable state has been proposed by (chandani, Takezoe et al. (xe2x80x9cJapanese Journal of Applied Physicsxe2x80x9d, vol. 27 (1988), pp. L729-). The anti-ferroelectric liquid crystal also provides a very quick response speed due to inversion switching based on spontaneous polarization similarly as in the ferroelectric liquid crystal.\nAs another type of the chiral smectic liquid crystal, there has been recently proposed a chiral smectic liquid crystal providing a V-character shaped response characteristic (voltage-transmittance characteristic) which is advantageous for gradational image display and is free from hysteresis (e.g., xe2x80x9cJapanese Journal of Applied Physicsxe2x80x9d, Vol. 36 (1997), pp. 3586-).\nFurther, an active matrix-type liquid crystal device using such a chiral smectic liquid crystal providing the V-shaped voltage-transmittance characteristic has also been proposed (JP-A 9-50049).\nIn recent years, the above-mentioned liquid crystal devices are required to be used for displaying motion (picture) images.\nIn the case where motion (picture) image are displayed by a liquid crystal panel, images to be displayed (still (picture) images) are changed for each frame period. In this case, if such a change in image is always recognized by a viewer, a transitional state of the image change is also consequently recognized, thus lowering image qualities of motion images. In order to solve the problem, a backlight (unit) is turned on only in a period wherein the still image display is completed in the liquid crystal panel.\nSuch a liquid crystal panel for displaying motion images is, however, accompanied with a problem of a hysteresis with respect to an alignment state of a liquid crystal used.\nSpecifically, such a hysteresis is a phenomenon that even when a prescribed voltage is applied for displaying a gradational (display) state of 50% in a frame period, the gradational state (level) of 50% cannot be realized by the influence of a gradational state in its preceding frame period.\nIn the conventional liquid crystal devices, in order to solve the above hysteresis phenomenon, a reset voltage has been applied in each frame period.\nMore specifically, in the conventional liquid crystal device, as shown by VR at (e) in FIG. 14, a fixed voltage (0 V in the figure) has been applied as a reset voltage. For this purpose, the liquid crystal device is required to additionally providing a reset circuit including a switching element 30 (for forcedly providing buffer circuit with a uniform potential), a terminal (for providing the buffer circuit with a potential corresponding to a reset potential), and a gate terminal 32 of the switching element 30 (for controlling the timing for supplying the potential corresponding to a reset potential to the buffer circuit via the switching element) as shown in FIG. 15, thus resulting in a complicated pixel circuit. FIG. 15 shows an embodiment of an equivalent circuit of the conventional liquid crystal device. Referring to FIG. 15, in addition to the reset circuit 30 (encoding the reset line 31 and the reset switching element 32), the conventional liquid crystal device includes a liquid crystal 1, a pair of electrodes 2a and 2b, a first switching element 3, a gate line 4, a signal line 5, a first storage (holding) capacitor 6, a first buffer circuit 7, a second switching element 8, a second buffer circuit 9, a common control line 10, a counter electrode potential 11, a common potential 12 and a second storage (holding) capacitor 13.\nWhen the conventional liquid crystal device as shown in FIG. 15 in driven for continuously displaying a black state in a certain pixel while setting a reset voltage of 0 V as shown at (b) in FIG. 16, a resultant voltage-transmittance characteristic (V-T characteristic) at the certain pixel as indicated by a curve connecting white squares (-xe2x96xa1-xe2x96xa1-) shown in FIG. 17 is different from a curve connecting black squares (-▪-▪-) shown in FIG. 17 indicating a V-T characteristic in the case of continuously displaying a white state at another pixel while setting a reset voltage of 0 V as shown at (a) in FIG. 16 (in this case, reset period =1 msec and writing period =8.33 msec are set), thus resulting in an occurrence of so-called image memory (burning or sticking) wherein a gradational image displayed on the liquid crystal panel as a whole is different from a gradational image to be displayed.\nFurther, in the case where the conventional liquid crystal device is driven for displaying full-color images according to a field-sequential driving scheme, a part of image data displayed in ia preceding frame period is displayed in a current frame period, thus resulting in an actually displayed color image which is different from a color image to be displayed originally, i.e., a poor color reproducibility.\nAn object of the present invention is to provide a driving method for a liquid crystal device capable of displaying appropriate gradational images while controlling a reset voltage for resetting a liquid crystal in a prescribed state without using an additional circuit (device) for exclusively applying a reset voltage.\nAnother object of the present invention is to provide a driving method for a liquid crystal device capable of suppressing an occurrence of image memory phenomenon.\nAccording to the present invention, there is provided a driving method for a liquid crystal device comprising a pair of electrodes and a liquid crystal disposed between the electrodes, said driving method comprising:\na sequence of voltage application operations each comprising application of a reset voltage to the liquid crystal for placing the liquid crystal in a reset state in a reset period and application of a data voltage to the liquid crystal for placing the liquid crystal in a desired gradational display state in a writing period subsequent to the reset period, wherein\neach reset voltage is set to provide a prescribed difference in voltage between said each reset voltage and a subsequent data voltage."}
{"text": "This invention relates to game apparatus permitting the retrieval of a ball from a random location within a location-concealing chamber.\nVarious designs of game apparatus are known, having a random nature. For example, U.S. Pat. No. 5,536,007 describes a cat toy having a ball located behind a screen with holes formed therethrough. The intention is that a cat may play with the toy, watching the ball through the holes, but there is no arrangement for retrieving the ball through the holes. In addition, U.S. Pat. No. 4,630,822 describes a game where a number of ball are randomly knocked into discrete apertures in a dish-shaped playing surface by a spinning top. The balls may be recovered from the various value apertures by a collecting device.\nAccording to the present invention, there is provided a game apparatus comprising: a lower surface; an upper surface over lying the lower surface and spaced therefrom so as to define a concealment chamber intermediate the upper and lower surfaces, the upper surface having a plurality of apertures each of which is provided with a closure means; a ball for positioning at a random location on the lower surface and being concealed within the concealment chamber; and a retrieval device connectable to the ball and which may be used to attempt retrieval of the ball from within the concealment chamber by inserting the retrieval device through a selected one of the plurality of apertures.\nIt will be appreciated that the ball may be disposed at a random location on the lower surface simply by shaking the apparatus and then the retrieval device is used to attempt to extract the ball from the chamber, by inserting that device into a randomly selected opening. Depending upon the configuration of the closure, it may be necessary to open that closure before inserting the device or the device may be inserted through the closurexe2x80x94but in either case, the design of the closure should be such that the location of the ball cannot be determined visually, from an external viewpoint.\nThe apparatus may be associated with a game-playing board carrying appropriate markings and in that case the ball, once retrieved from the apparatus, may be employed in further game play, on that board. In this case, the apparatus may be configured so as to be connectable to and easily removable from the game-playing board.\nPreferably the ball is made from a ferro-magnetic material such as steel, and the retrieval device has a magnetic end for releasable connection to the ball. Such a retrieval device is conveniently rod shaped.\nWhilst the apparatus could be used separately from a game board, the apparatus may advantageously be placed on a board and the board moved from side to side by the players. These movements enable a ball placed within the apparatus to roll to an unknown random location within the chamber. The retrieval of the ball from its concealed position can then be attempted by the insertion of the retrieval device through a randomly-selected aperture in the upper surface. The apparatus thus allows the players in a game the opportunity of locating and retrieving a ball, by chance.\nIt is desirable that the lower surface defines a number of specific ball locating positions, such that in use the ball will normally remain in one position, but if sufficient force is applied to the device, the ball will be dislodged from this one position and move within the intermediate concealment chamber, randomly to relocate at any one of the ball locating positions. For example, the specific ball locating positions may take the form of peaks and troughs or recesses provided on the lower surface, each such ball locating position being disposed directly beneath an aperture in the upper surface.\nThe closure may take any of a number of forms, but preferably each closure comprises a pair of shutters hingedly connected to the upper surface, and movable between respective closed positions where they are parallel to the upper surface and cover the aperture and open positions where they permit the removal of a concealed ball. Alternatively, each closure may comprise a removable block that covers the aperture but is completely removable therefrom to permit an attempt at removal of a concealed ball.\nIt is preferred that the closure should conceal the location of a ball, whilst permitting the entry of the retrieval device and removal of the ball with the least amount of actions. To this end a closure that allows entry of the ball retrieval device through a hole therein and then removal of the ball without any separate opening or closing action is preferred.\nSuch a hole should be suitably sized to permit the passage of the retrieval device therethrough but prevent a concealed ball from being visible from outside the device. Therefore, an attempt at removal of the bail is a simple one-handed task wherein the retrieval device is inserted through the chosen hole without hindrance from the closure; if the ball is present beneath the aperture the retrieval device makes contact with the ball to permit its removal. This is achieved by pulling the retrieval device out of the aperture, the passage of the ball through the aperture serving to open the closure.\nIn the alternative, a closure that is not openable automatically by the retrieval device or the ball may be employed. For example, each closure may comprise a single lid hingedly connected to the upper surface and movable between a closed position covering the aperture and an open position where entry of the retrieval device and removal of a ball is permitted. Alternatively, each closure may comprise a sliding door connected to the upper surface and movable between open and closed positions. To facilitate the movement of such a sliding door, runners may be arranged on the upper surface of the device, for guiding and retaining all of the doors.\nAs the ball is likely to be removed during game play, a side wall of the apparatus is preferably provided with an access hatch having an openable cover, to permit the reintroduction of the ball into the intermediate chamber. In addition, the game apparatus is preferably provided with a recess on the underside of the lower surface to allow the apparatus to be positioned and secured on a game board.\nIn use, the apparatus is shaken randomly to position the ball within the intermediate chamber. A retrieval device is inserted through an aperture in the upper surface, either after opening the associated closure or through the closure, depending on the type of closure used. If the ball, or one of a number of balls, contained within the chamber is found beneath the aperture, the retrieval device connects to the ball and removal of the retrieval device extracts the ball from the chamber."}
{"text": "The present invention relates to materials and methods of activating T-cells with specificity for particular antigenic peptides, the use of activated T-cells in vivo for the treatment of a variety of disease conditions, and compositions appropriate for these uses.\nThe efficiency with which the immune system cures or protects individuals from infectious disease has always been intriguing to scientists, as it has been believed that it might be possible to activate the immune system to combat other types of diseases. Such diseases include cancer, AIDS, hepatitis and infectious disease in immunosuppressed patients. While various procedures involving the use of antibodies have been applied in those types of diseases, few if any successful attempts using cytotoxic T-cells have been recorded. Theoretically, cytotoxic T-cells would be the preferable means of treating the types of disease noted above. However, no procedures have been available to specifically activate cytotoxic T-cells.\nCytotoxic T-cells, or CD8xe2x80x2 cells (i.e., cells expressing the molecule CD8) as they are presently known, represent the main line of defense against viral infections. CD8xe2x80x2 lymphocytes specifically recognize and kill cells which are infected by a virus. Thus, the cost of eliminating a viral infection is the accompanying loss of the infected cells. The T-cell receptors on the surface of CD8+ cells cannot recognize foreign antigens directly. In contrast to antibodies, antigen must first be presented to the receptors.\nThe presentation of antigen to CD8+ T-cells is accomplished by major histocompatibility complex (MHC) molecules of the Class I type. The major histocompatibility complex (MHC) refers to a large genetic locus encoding an extensive family of glycoproteins which play an important role in the immune response. The MHC genes, which are also referred to as the HLA (human leucocyte antigen) complex, are located on chromosome 6 in humans. The molecules encoded by MHC genes are present on cell surfaces and are largely responsible for recognition of tissue transplants as xe2x80x9cnon-selfxe2x80x9d. Thus, membrane-bound MHC molecules are intimately involved in recognition of antigens by T-cells.\nMHC products are grouped into three major classes, referred to as I, II, and III. T-cells that serve mainly as helper cells express CD4 and primarily interact with Class II molecules, whereas CD8-expressing cells, which mostly represent cytotoxic effector cells, interact with Class I molecules.\nClass I molecules are membrane glycoproteins with the ability to bind peptides derived primarily from intracellular degradation of endogenous proteins. Complexes of MHC molecules with peptides derived from viral, bacterial and other foreign proteins comprise the ligand that triggers the antigen responsiveness of T-cells. In contrast, complexes of MHC molecules with peptides derived from normal cellular products play a role in xe2x80x9cteachingxe2x80x9d the T-cells to tolerate self-peptides, in the thymus. Class I molecules do not present entire, intact antigens; rather, they present peptide fragments thereof, xe2x80x9cloadedxe2x80x9d onto their xe2x80x9cpeptide binding groovexe2x80x9d.\nFor many years, immunologists have hoped to raise specific cytotoxic cells targeting viruses, retroviruses and cancer cells. While targeting against viral diseases in general may be accomplished in vivo by vaccination with live or attenuated vaccines, no similar success has been achieved with retroviruses or with cancer cells. Moreover, the vaccine approach has not had the desired efficacy in immunosuppressed patients. At least one researcher has taken the rather non-specific approach of xe2x80x9cboostingxe2x80x9d existing CD8xe2x80x2 cells by incubating them in vitro with IL-2, a growth factor for T-cells. However, this protocol (known as LAK cell therapy) will only allow the expansion of those CD8xe2x80x2 cells which are already activated. As the immune system is always active for one reason or another, most of the IL-2 stimulated cells will be irrelevant for the purpose of combatting the disease. In fact, it has not been documented that this type of therapy activates any cells with the desired specificity. Thus, the benefits of LAK cell therapy are controversial at best, and the side effects are typically so severe that many studies have been discontinued.\nSeveral novel molecules which appear to be involved in the peptide loading process have recently been identified. It has also been noted that Class I molecules without bound peptide (i.e., xe2x80x9cemptyxe2x80x9d molecules) can be produced under certain restrictive circumstances. These xe2x80x9cemptyxe2x80x9d molecules are often unable to reach the cell surface, however, as Class I molecules without bound peptide are very thermolabile. Thus, the xe2x80x9cemptyxe2x80x9d Class I molecules disassemble during their transport from the interior of the cell to the cell surface.\nThe presentation of Class I MHC molecules bound to peptide alone has generally ineffective in activating CD8xe2x80x2 cells. In nature, the CD8+ cells are activated by antigen-presenting cells which present not only a peptide-bound Class I MHC molecule, but also a costimulatory molecule. Such costimulatory molecules include B7 which is now recognized to be two subgroups designated as B7.1 and B7.2. It has also been found that cell adhesion molecules such as integrins assist in this process.\nWhen the CD8+ T-cell interacts with an antigen-presenting cell having the peptide bound by a Class I MHC and costimulatory molecule, the CD8+ T-cell is activated to proliferate and becomes an armed effector T-cell. See, generally, Janeway and Travers, Immunobiology, published by Current Biology Limited, London (1994), incorporated by reference.\nAccordingly, what is needed is a means to activate T-cells so that they proliferate and become cytotoxic. It would be useful if the activation could be done in vitro and the activated cytotoxic T-cells reintroduced into the patient. It would also be desirable if the activation could be done by a synthetic antigen-presenting matrix comprised of a material such as cells which not only presents the selected peptide, but also presents other costimulatory factors which increase the effectiveness of the activation.\nIt would also be advantageous if it was possible to select the peptide so that substantially only those CD8+ cells cytotoxic to cells presenting that peptide would be activated.\nThe present invention relates to a synthetic antigen-presenting system for presenting an MHC molecule complexed to a peptide and an assisting molecule to a T-cell to activate the T-cell.\nIn one embodiment, the system relates to a synthetic antigen-presenting matrix having a support and at least the extracellular portion of a Class I MHC molecule capable of binding to a selected peptide operably linked to the support. The matrix also includes an assisting molecule operably linked to the support. The assisting molecule acts on a receptor on the CD8+ T-cell. The MHC and assisting molecules are present in sufficient numbers to activate a population of T-cell lymphocytes against the peptide when the peptide is bound to the extracellular portion of the MHC molecule.\nIt has been found that an antigen-presenting matrix having both an MHC molecule or a portion of a MHC molecule together with an assisting molecule, provides a synergistic reaction in activating T-cell lymphocytes against the peptide. Examples of assisting molecules are costimulatory molecules such as B7.1 and B7.2 or adhesion molecules such as ICAM-1 and LFA-3. The extracellular portion of such costimulatory molecules can also be used. Another type of costimulatory molecule is one that reacts with the CD28 molecule such as anti-CD28 antibodies or functional portions thereof, e.g. Fab portions.\nIt has been found that a specifically effective synergistic reaction results from an antigen-presenting matrix having MHC molecules bound with a peptide, a costimulatory molecule, and an adhesion molecule. In particular, a highly effective synergistic generation of cytotoxic T-cell activity results from the combination of 27.1 and ICAM-1.\nThe support used for the matrix can take several different forms. Examples for the support include solid support such as metals or plastics, porous materials such as resin or modified cellulose columns, microbeads, microtiter plates, red blood cells and liposomes.\nAnother type of support is a cell fragment, such as a cell membrane fragment or an entire cell. In this embodiment, the matrix is actually cells which have been transfected to present MHC molecules and assisting molecules on the cell surface to create an antigen-presenting cell (APC). This is done by producing a cell line containing at least one expressible Class I MHC nucleotide sequence for the MHC heavy chain, preferably a cDNA sequence, operably linked to a first promoter and an expressible xcex2-2 microglobulin nucleotide sequence operably linked to a second promoter. The MHC heavy chain and the xcex2-2 microglobulin associate together form the MHC molecule which binds to the peptide. The MHC protein binds with the antigenic peptide and presents it on the surface of the cell. The cell also includes a gene for a nucleotide sequence of an assisting molecule operably linked to a third promoter. The assisting molecule is also presented on the surface of the cell. These molecules are presented on the surface of the cell in sufficient numbers to activate a population of T-cell lymphocytes against the peptide when the peptide is bound to the complexes. Other molecules on the surface of a cell or cell fragment such as carbohydrate moieties may also provide some costimulation to the T-cells.\nThe cell line is synthetic in that at least one of the genes described above is not naturally present in the cells from which the cell line is derived. It is preferable to use a poikilotherm cell line because MHC molecules are thermolabile. A range of species are useful for this purpose. See, for example, U.S. Pat. No. 5,314,813 to Petersen et al. which discusses numerous species for this use and is incorporated by reference. It is preferred to use eukaryotic cells and insect cells in particular.\nIn one embodiment, it is particularly preferred to have at least two assisting molecules, one being a costimulatory molecule and the other being an adhesion molecule. It has been found that this combination has a synergistic effect, giving even greater T-cell activation than either of the individual molecules combined. It has also been found to be advantageous and preferable to have at least one of the transfected genes under control of an inducible promoter.\nUsing the present invention, it is possible to introduce the peptide to the cell while it is producing MHC molecules and allow the peptide to bind the MHC molecules while they are still within the cell. Alternatively, the MHC molecules can be expressed as empty molecules on the cell surface and the peptide introduced to the cells after the molecules are expressed on the cell surface. In this latter procedure, the use of poikilotherm cells is particularly advantageous because empty MHC molecules, those not yet complexed or bound with peptides, are thermolabile.\nClass I MHC molecules have been expressed in insect cells such as Drosophila melanogaster (fruit fly) cells. Since Drosophila does not have all the components of a mammalian immune system, the various proteins involved in the peptide loading machinery should be absent from such cells. The lack of peptide loading machinery allows the Class I molecules to be expressed as empty molecules at the cell surface.\nAnother advantage of using insect cells such as the Drosophila system is that Drosophila cells prefer a temperature of 28xc2x0 C. rather than 37xc2x0 C. This fact is very important, because empty Class I molecules are thermolabile and tend to disintegrate at 37xc2x0 C. By incubating the Class I-expressing Drosophila cells with peptides that can bind to the Class I molecule, it is possible to get virtually every Class I molecule to contain one and the same peptide. The cells are accordingly very different from mammalian cells, where the Class I molecules contain many different types of peptides, most of which are derived from our own, innocuous cellular proteins.\nThe present invention also relates to methods for producing activated CD8+ cells in vitro. One method comprises contacting, in vitro, CD8+ cells with one of the antigen-presenting matrices described above for a time period sufficient to activate, in an antigen-specific manner, the CD8+ cells. The method may further comprise (1) separating the activated CD8+ cells from the antigen-presenting matrix; (2) suspending the activated CD8+ cells in an acceptable carrier or excipient; and (3) administering the suspension to an individual in need of treatment. The antigens may comprise native or undegraded proteins or polypeptides, or they may comprise antigenic polypeptides which have been cleaved into peptide fragments comprising at least 8 amino acid residues prior to incubation with the human Class I MHC molecules.\nIn another variation, the invention relates to methods treating conditions in patients and specifically killing target cells in a human patient. The method comprises (1) obtaining a fluid sample containing resting or naive CD8+ cells from the patient; (2) contacting, in vitro, the CD8+ cells with an antigen-presenting matrix for a time period sufficient to activate, in an antigen-specific manner, the CD8+ cells; and (3) administering the activated CD8+ cells to the patient. For example, the use of tumor specific peptides allows for the treatment of tumor related diseases by producing cytotoxic activated CD8+ T-cells. The invention also relates to the method of treating a medical condition by administration of an antigen-presenting matrix in a suitable suspension. In various embodiments the condition may comprise cancer, tumors, neoplasia, viral or retroviral infection, autoimmune or autoimmune-type conditions. In one embodiment, the method of administering the matrix comprises intravenous injection."}
{"text": "Commonly available stethoscopes provide a monaural sound in that the sound is received in a single sound accumulator such as bell type accumulators or diaphragm type accumulators from which it is \"mechanically\" provided to ear pieces via ear tubes. A person using such a conventional, monophonic stethoscope will not be able to perceive any directionality of spatiality of the sounds heard. This is because the sound is transmitted to both ear tubes from the one sound accumulator, which acts as the transducer of the sound and therefore becomes a single, intermediate source, in which the sounds generated by the true sources are superimposed.\nThe perception of directionality and spatiality of the sounds in stethoscopes would allow a differentiation between sounds. E.g cardiac murmurs which are the result of a turbulence of the blood flow in the heart and thus the origins of sounds associated with cardiac murmurs are widely spread whereas the closing and opening clicks of cardiac valves originate from a point source.\nTo enable perception of directionality and spatiality of sound obtained from a stethoscope, each ear must be provided with sounds transduced by separate sound accumulators and consequently two sound channels must be provided to realise a stereophonic stethoscope.\nStethoscopes providing a stereo signal mechanically have been produced and tested, e.g. reference is made to S. Kazama, \"A new Stereophonic Stethoscope\", Jpn. Heart J, Vol. 31/6, November 1990.\nOn the other hand, the use of the conventional, mechanical stethoscope principle has limitations in that sounds can only be heard live in the sense that typically only one person is listening to the sounds and bases his or her diagnosis on his perception of the sounds. To obtain a second opinion of another person, this person must also perform an examination of the patient with the stethoscope therefore the patient is required for a longer time and discomfort may be imposed onto the patient. Therefore, the use of electronic stethoscopes can enable an easy and direct way to record the sounds in one examination which can be replayed to other persons at any time and without further examination of the patient. Further, the electronic audio signal can be more readily made available to more than one person during the examination by transducing the audio signal via speakers. An electronic stethoscope in which a microphone is connected to one precordial stethoscope probe or esophageal probe to monitor a monophonic sound accumulated in the bell or probe electronically via a high quality amplifier is described by J. H. Philip et al in \"An electronic stethoscope is judged better than conventional stethoscopes for anaesthesia monitoring\", Journal of Clinical Monitoring, Vol. 2/3, July 1986. Previous electronic stethoscopes have employed electronic frequency filtering to selectively amplify sound in the range 37.5-1000 Hz where most heart and lung sounds occur 1, e.g. a bandpass frequency of 0.7-5000 Hz with selectable bass and treble boost/cut circuitry. Other electronic stethoscopes use an upward frequency shift of the entire audio band, aiming to allow improved detection of low frequency sounds. However, the electronic manipulation involved makes for complex and expensive circuitry, as well as increased size, particularly if two units are to be employed for stereophonic stethoscopes. It also requires a significant degree of retraining as the familiar heart and lung sounds may be quite altered."}
{"text": "Mobile telephones and similar communication devices are rapidly expanding in use and function. Such devices will soon provide Internet access, personal information management, facsimile, messaging, in addition to telephone communication. To accomplish this there is a need to provide keyboards compatible with the more complex applications to which the mobile device will be adapted. Full function keyboards, such as the standard QWERTY typing array of keys and buttons, are difficult to provide while maintaining the compact size required in the mobile device. Such devices on the market today are cumbersome and often require a separate belt pouch for carrying the mobile device on the person of the user. In addition it is necessary to coordinate on screen displays for adaptation to the variety of functions.\nIt is a purpose of this invention to provide a simple and inexpensive means of providing a full function keyboard to accommodate the burgeoning applications to which a mobile communication device is adapted.\nA mobile communication device is constructed having a body in which is mounted a viewing screen for displaying user interface communications. A first panel is fixed to the body and includes one half of a full function keyboard and a back cover. A second panel is pivotally mounted on the body for rotation between two positions, namely, a closed position and an open position. The second panel is constructed with inner and outer surfaces located on opposite sides of the panel which are positioned such that the inner surface faces the keyboard of the first panel and the outer surface is exposed in the closed position. In the open position, the second panel is pivoted approximately 180xc2x0 so that the outer surface faces away from the user of the mobile device while the inner surface is exposed.\nTo enable the device to be operated as a communication device a communication keypad is constructed in the outer surface of the second panel. To provide the second half of the full function keyboard, the inner surface contains a keyboard which is operatively associated with the keyboard of the first panel in the open position. The screen remains exposed in each position of the second panel. The display which appears on the screen is oriented in a manner consistent with the position of panels.\nIn another embodiment of this invention, a switch is operatively associated with the moveable second panel to send a signal to the microprocessor controller of the device indicative of the position of the second panel. This information is used by the controller to adjust the function of keys to be consistent with the application for which the second panel is positioned. In particular an array of soft keys is mounted on the body of the display for convenient use. These soft keys have different functions depending on the position of the second panel.\nIn a preferred embodiment of this invention an array of three soft keys are positioned about the periphery of the user so that there are always two keys available at the bottom of the display and a third key available at the upper portion of the display. In this arrangement the key pattern is substantially the same irrespective of orientation of the device. The orientation will be horizontal or vertical depending on the position of the second panel. It is a purpose of this invention to provide substantially the same key functions for the key in a particular relative position in both modes of operation. The upper key, right or left, will be the power key and the lower right and left keys will each operate as a soft key to open quick selection profiles. Each key may be assigned a different menu for different key functions.\nConsidering first the use of the device with the second panel closed, for example as a mobile telephone, soft keys are arranged at the lower left, the lower right and the upper left of the display screen. In this position, the upper left hand key will operate as the power on/off key. The lower right and left keys operate as a soft key to open quick selection profiles and menus. For an example, one function of the lower keys that may be selected would be as up and down and left and right scroll keys. Other alternate functions may be devised depending on the applications served.\nWhen the second panel is opened, the assigned function of the lower soft keys is shifted clockwise or counter clockwise by the controller depending on the use orientation of the device with the second panel open. The functions assigned to the newly positioned lower left and right keys will be the same as the keys positioned at these locations in the closed position. The upper key moves to the upper right and will be the power key. The device controller adjusts the assignment of key functions available at a soft key to provide a similar set of functions at substantially the same relative key location in the open position of the second panel, as in the closed position."}
{"text": "1. Field of the Invention\nThe present invention relates to image processing performed to protect a copyright for an image, and the like.\n2. Related Background Art\nRecently, the standardization of motion image encoding techniques has been advanced rapidly. And encoding standards, such as MPEG-1, MPEG-2, H.261 and H.263 (these encoding systems are specifically accepted as international standards by ISO (International Organization for Standardization) and ITU (International Telecommunication Union)), have now been established. Further, work is progressing on MPEG-4, a new international standard that is intended to be a general-purpose, next-generation multimedia encoding standard that can be used in many fields, such as computers, broadcasting and communication. Image digitization has also been advanced by the encoding technique, and the storage and communication of motion image data have been markedly improved.\nFor example, with MPEG-4, a motion image scene at a specific moment is regarded as an aggregate of objects, such as a “person”, a “car”, a “background”, “music” and “video”, which are the components of the scene, and an appropriate encoding process is performed for each object to efficiently encode all the motion image data.\nAccording to a multi-object encoding technique such as MPEG-4, a frame for a scene is expressed by combining individual objects that have been encoded and compressed. Therefore, object data obtained from specific MPEG-4 data can be re-used with different MPEG-4 data, so that another frame can be easily expressed.\nAs the above described digital encoding standards have spread, the contents business has strongly expressed the problem they represent for the provision of copyright protection. That is, if the standards do not satisfactorily ensure adequate copyright protection, safe and superior contents can not be provided.\nFor MPEG-4 this is a very serious problem. This is so because, as is described above, with MPEG-4 it is possible to encode and to operate/process data for individual objects. Therefore, since the use of the multi-object encoding system MPEG-4 permits the seamless reuse of individual objects, it is natural for requests to be made that protection be provided for the intellectual property represented by an object. And if the target is a person, it is natural that protection be sought for the right of portrait.\nTo protect a copyright for a portion of a motion image, a method is available for temporarily halting the decoding of the motion image to prevent the reproduction of a part of it. According to this method, while taking into account the right of portrait and the copyright for the motion image, the decoding of a portion related to them is halted until that portion has been passed, at which time the decoding is resumed.\nWith this method, however, the following problem is encountered.\nThe encoding system for using interframe correlation is generally employed to encode motion images. In the previously mentioned, well-known encoding systems H.261 and H.263, and MPEG-1, 2 and 4, basically, the preceding frame, or both the preceding and the succeeding frames, are referred to along the time axis, and motion compensation is performed to encode motion images.\nData reproduction provided by H.261 or H.263 is shown in FIG. 1. I* denotes a frame for intraframe encoding, and P* denotes a frame for interframe encoding. In FIG. 1, Time is the direction in which time elapses, and in Security, a period during which the decoding process is halted is represented as a black block. Code represents the arrangement of screens in the order in which encoded, and Display represents the arrangement of screens in the order in which displayed.\nSuppose that the decoding of the encoded data for frames P4 to P9 is halted to apply the protection (Security) provided for an intellectual property (e.g., a copyright). The decoding of motion images is halted at data P3, and images are not again displayed until the decoding is resumed. Since the writing of encoded data to a buffer is also halted at the same time, the encoded data from P4 to P9 are abandoned. Therefore, when the image protection time has ended and the decoding is resumed, data P9, which originally should have been referred to by data P10, have already been abandoned, so that data P10 and the following data can not be normally decoded. As a result, until frame I1, for which intraframe encoding is performed, is decoded, in an interval extending from P10 to P13, the images in the frames are deteriorated, or the decoding process is halted.\nMPEG-1, 2, and 4 also employ bidirectional motion compensation, whereby motion along the time axis is predicted from preceding and succeeding frames.\nThe data reproduction provided by MPEG-1, 2, and 4 is shown in FIG. 2. B* denotes frames for bidirectional (two-way) interframe encoding. In FIG. 2, Time is the direction in which time elapses, and in Security, a period during which the decoding process is halted is represented as a black block. Code represents the arrangement of screens in the order in which encoding is performed, and Display represents the arrangement of screens in the order in which of displayed. For the frames B, the order in which pictures are encoded differs from the order for the original pictures.\nSuppose that the decoding for encoded data from B10 to P3 is halted to protect an intellectual property (e.g., a copyright). Since motion compensation is not correctly performed, as in FIG. 1, the images in frames B30 and B31, or B40 and B41, are not correctly decoded, and the images in the frames are deteriorated, or the decoding process is halted."}
{"text": "Problems facing the use of onboard electronics in vehicles include the lack of diagnostic ability by service personnel and No Fault. Found (NFF) warranty costs. When a controller is returned to the manufacturer, it must be assumed to be defective until proven otherwise. Service and Warranty departments need time to put a returned controller on a manufacturing tester and perform other diagnostic tasks to determine either the root cause of the problem or to determine whether it is a NFF controller, i.e., to determine that no problem could be found. Some controllers that are returned as NFF may indeed have a physical problem, but the problem may not manifest itself unless it is embedded in a vehicle and is being subjected to vibrations, temperature extremes, or other operating conditions."}
{"text": "Convolution is commonly performed on signals in many contexts, including the fields of sound, still image, video, lithography, and radio (radar) signal processing. Typically, the signals to be convolved are pattern signals. Each of the expressions “pattern” and “pattern signal” is used herein in a broad sense to denote a one-dimensional sequence or two-dimensional (or higher dimensional) array of data words (which can be, but need not be pixels). Typically, the data words comprise binary bits, and the convolution is performed in discrete fashion on the binary bits using software, digital signal processing circuitry, custom hardware, or FPGA systems (field programmable gate array based computing systems).\nThe term “data” herein denotes one or more signals indicative of data, and the expression “data word” herein denotes one or more signals indicative of a data word.\nThe motivations for implementing convolution rapidly, even when processing data indicative of very large patterns, are myriad. The present invention was motivated by the need for proximity correction in the field of lithography. In such problems, one attempts a two-dimensional convolution between data indicative of a large pattern “p” (where the pattern is a pixel array) and a diffusion kernel “d”. Often the kernel “d” is a Gaussian or a superposition of Gaussians, or is otherwise a smooth kernel. More specifically, the present invention grew out of attempts to establish a suitable “O(NN′)” algorithm (an algorithm requiring not more than on the order of NN′ multiplications and additions) for convolving a two-dimensional pattern comprising NN′ pixels, where each of N and N′ is very large) with a Gaussian kernel (or other smooth kernel) such that the convolution is exact or very close to exact.\nThe objective in performing proximity correction (in the field of lithography) is to generate a “raw” optical signal (or “raw” electron beam signal) which can be input to a set of reflective or refractive optics (or electron beam optics), in order to cause the output of the optics to produce a desired pattern on a mask or wafer. To determine the characteristics of a raw optical signal (or raw electron beam signal) that are needed to produce the desired pattern on the mask or wafer, a deconvolution operation is typically performed on a very large array of pixels (which determine a pattern “p”) in order to correct for the well known proximity problem. In the case of electron beam lithography, the proximity problem results from electron scattering in the substrate (mask or wafer) being written. Such scattering exposes broadened areas on the substrate to electrons (i.e., an area surrounding each pixel to be written in addition to the pixel itself), with the scattering effectively broadening the electron beam beyond the beam diameter with which the beam is incident on the substrate.\nIn nearly all proximity correction schemes, such a deconvolution operation includes at least one convolution step. Accordingly, in performing typical proximity correction, a very large array of pixels (determining a pattern “p”) must be convolved with a diffusion kernel. Although such a convolution is typically performed on a pattern comprising a very large array of binary pixels, this restriction is not essential in the following discussion and is not essential to implementation of the invention. Indeed, the invention can implement convolution on data indicative of any pattern “p” with a smooth convolution kernel “d” having characteristics to be described below.\nFor data indicative of a pattern “p” and a convolution kernel “d” we consider the cyclic convolution:                     (                  d          ⁢                      ×            C                    ⁢          p                )            n        =                  ∑                              i            +            j                    ≡                      n            ⁢                          (                              mod                ⁢                                                                   ⁢                N                            )                                          ⁢                        d          i                ⁢                  p          j                      ,where ×C denotes that the convolution operator has cyclic character, and an acyclic convolution which differs only in the indicial constraint and range:                     (                  d          ⁢                      ×            A                    ⁢          p                )            n        =                  ∑                              i            +            j                    =          n                    ⁢                        d          i                ⁢                  p          j                      ,where “×A” denotes that convolution operator has acyclic character.\nFor simplicity, we restrict much of the discussion herein to one-dimensional cases (in which the pattern p is an ordered set of N data values and the kernel is an ordered set of M values). Despite this, it should be appreciated that in typical embodiments of the invention, the pattern is two-dimensional (a two-dimensional array pjk of data values determines the pattern) and the summation defining the convolution (a summation which corresponds to either one of the summations set forth in the previous paragraph) is over index k as well as index j of the array pjk. In the case of a two-dimensional pattern p determined by an N by N′ array of data values, the indices n, i, j and domain lengths N in the formulae set forth in the previous paragraph are 2-vectors.\nIn one-dimensional cases, the result of the cyclic convolution has length N (it comprises N data values), and the result of the acyclic convolution has length M+N−1.\nIt is standard that a cyclic convolution d×Cp can be cast as an equivalent matrix-vector product:d×Cp≡Dp,where D is the circulant matrix of d (hereinafter the “circulant” of d), whose 1-dimensional form is defined (assuming that N is greater than 3) as:   D  =            (                                                  d              0                                                          d              1                                                          d              2                                                          d              3                                            ⋯                                              d                              N                -                1                                                                                        d                              N                -                1                                                                        d              0                                                          d              1                                                          d              2                                            ⋯                                              d                              N                -                2                                                                          ⋮                                                                                                                                                                                                                                                                            ⋮                                                              d              1                                                          d              2                                                          d              3                                                          d              4                                            ⋯                                              d              0                                          )        .  Therefore, conventional methods for cyclic convolution can be cast in the language of matrix algebra. Acyclic convolution can be obtained with similar matrix manipulations.\nFor the sake of simplicity, we will use the symbol × hereinbelow to denote convolution having either acyclic or cyclic character. In most of the discussion the symbol will refer to convolution having acyclic character. Those of ordinary skill in the art will recognize that, given a specified acyclic convolution, a corresponding cyclic convolution can be implemented by slight modification of the parameters (e.g., the boundary conditions and definition of the circulant of the kernel) that determine the acyclic convolution.\nAbove-referenced U.S. patent application Ser. No. 09/480,908 discloses a fast convolution method whose central idea (in one-dimensional embodiments) is to approximate a smooth kernel d by a polynomial spline kernel ƒ (where ƒ is a spline function ƒ(x) which is piecewise a polynomial of degree δ with L pieces f1(x)), and then to use appropriate operators that annihilate (or flatten) each polynomial of given degree (in a manner to be explained) to calculate the convolution of ƒ and p quickly. In some embodiments, the smooth kernel d is approximated by a spline kernel ƒ which is not a polynomial spline kernel, but which consists of L pieces defined over adjacent segments of its domain (in typical two-dimensional cases, the latter spline kernel is a radially symmetric function whose domain is some continuous or discrete set of values of the radial parameter). Though “spline” convolution as described in U.S. application Ser. No. 09/480,908 has features reminiscent of conventional wavelet schemes and is an O(N) algorithm (as are wavelet schemes), an advantage of “spline” convolution is that it can be performed (on data indicative of a pattern p consisting of N data values) with cN arithmetic operations (multiplications and additions), whereas conventional wavelet convolution on the same data would require bN arithmetic operations, where the factor “b” is typically (i.e., with typical error analysis) significantly larger than the factor “c.” In other words, the implied big-O constant for the spline convolution is significantly smaller than the typical such constant for conventional wavelet convolution.\nSpline convolution, as described in U.S. application Ser. No. 09/480,908, is a method for performing cyclic or acyclic convolution of a pattern “p” (i.e., data indicative of a pattern “p”) with a smooth diffusion kernel d, to generate data indicative of the convolution result r=Dp, where D is the circulant of d. The pattern “p” can be one-dimensional in the sense that it is determined by a continuous (or discrete) one-dimensional domain of data values (e.g., pixels), or it can be two-dimensional in the sense that it is determined by a continuous two-dimensional domain of data values (or a two-dimensional array of discrete data values), or p can have dimension greater than two. In typical discrete implementations, the pattern p is one-dimensional in the sense that it is determined by a discrete, ordered set of data values (e.g., pixels) pi, where i varies from 0 to N−1 (where N is the signal length), or it is two-dimensional in the sense that it is determined by an array of data values pij, where i varies from 0 to N−1 and j varies from 0 to N′−1, or it has dimension greater than two (it is determined by a three or higher-dimensional set of data values). Typically, the kernel d is determined by an array of data values dij, where i varies from 0 to N−1 and j varies from 0 to N′−1 (but the kernel d can alternatively be determined by a discrete set of data values d0 through dN−1).\nIn some embodiments described in U.S. application Ser. No. 09/480,908, the convolution Dp is accomplished by performing the steps of:\n(a) approximating the kernel d by a polynomial spline kernel ƒ (unless the kernel d is itself a polynomial spline kernel, in which case d=ƒ and step (a) is omitted);\n(b) calculating q=Bp=Δδ+1Fp, where F is the circulant of kernel ƒ, and Δδ+1 is an annihilation operator (whose form generally depends on the degree δ of the polynomial segments of ƒ) which operates on circulant F in such a manner that Δδ+1F=B is sparse; and\n(c) backsolving Δδ+1r=q to determiner=Fp.\nIn cases in which the kernel d is itself a polynomial spline kernel (so that d=ƒ and F=D), the method yields an exact result (r=Dp). Otherwise, the error inherent in the method is (ƒ−d)×p, where × denotes convolution, and thus the error is bounded easily.\nIn one-dimensional cases (in which the pattern to be convolved is a one-dimensional pattern of length N), Δδ+1 has the form of the N×N circulant matrix defined as follows:       Δ          δ      +      1        =      [                                        +                          (                                                                                          δ                      +                      1                                                                                                            0                                                              )                                                            -                          (                                                                                          δ                      +                      1                                                                                                            1                                                              )                                                            +                          (                                                                                          δ                      +                      1                                                                                                            2                                                              )                                                            -                          (                                                                                          δ                      +                      1                                                                                                            3                                                              )                                                ⋯                          0                                      0                                      +                          (                                                                                          δ                      +                      1                                                                                                            0                                                              )                                                            -                          (                                                                                          δ                      +                      1                                                                                                            1                                                              )                                                            +                          (                                                                                          δ                      +                      1                                                                                                            2                                                              )                                                ⋯                          0                                      ⋮                                                                                                                                                                                                                              ⋮                                                  -                          (                                                                                          δ                      +                      1                                                                                                            1                                                              )                                                            +                          (                                                                                          δ                      +                      1                                                                                                            2                                                              )                                                ⋯                                                                           0                                      +                          (                                                                                          δ                      +                      1                                                                                                            0                                                              )                                            ]  in which each entry is a binomial coefficient, and δ is the maximum degree of the spline segments of spline kernel ƒ. For example, δ=2 where the spline kernel ƒ comprises quadratic segments. In two- or higher-dimensional cases, the annihilation operators can be defined as       Δ    =                            ∂                      x            1                                n            1                          ⁢                              ∂                          x              2                                      n              2                                ⁢                                           ⁢          …                    ⁢                           ⁢              ∂                  x          d                          n          d                      ,where ∂xhnh is the n-th partial derivative with respect to the h-th of d coordinates. For example, the Laplacian       ∇    2    ⁢      =                            ∂                      x            1                    2                ⁢                  +                                           ⁢          …                    ⁢                           +              ∂                  x          d                2            will annihilate piecewise-planar functions of d-dimensional arguments ƒ=ƒ(x1, x2, . . . xd).\nIn the one-dimensional case, the end points of each segment (the “pivot points”) of spline kernel ƒ may be consecutive elements di and di+1 of kernel d, and step (a) can be implemented by performing curve fitting to select each segment of the spline kernel as one which adequately matches a corresponding segment of the kernel d. In some implementations, appropriate boundary conditions are satisfied at each pivot point, such as by derivative-matching or satisfying some other smoothness criterion at the pivot points.\nIn some implementations described in application Ser. No. 09/480,908, step (c) includes a preliminary “ignition” step in which a small number of the lowest components of r=Fp are computed by exact multiplication of p by a few rows of F, and then a step of determining the rest of the components of r using a natural recurrence relation determined by the spline kernel and the operator Δδ+1. For example, in the one-dimensional case, the lowest components of r are r0, r1, . . . , rδ, where “δ” is the maximum degree of the spline segments of spline kernel ƒ (for example r0, r1, and r2 where the spline kernel comprises quadratic segments), and these (δ+1) components are determined by exact multiplication of p by (δ+1) rows of F. The (δ+1) components can alternatively be determined in other ways. Then, the rest of the components “rδ” are determined using a natural recurrence relation determined by the operator Δδ+1. The “ignition” operation which generates the components r0, r1, . . . , rδ, can be accomplished with O(N) computations. The recurrence relation calculation can also be accomplished with O(N) computations.\nIn other embodiments, the method disclosed in U.S. application Ser. No. 09/480,908 for performing the convolution r=Dp includes the steps of:\n(a) approximating the kernel d by a polynomial spline kernel ƒ (unless the kernel d is itself a polynomial spline kernel, in which case d=ƒ and step (a) is omitted);\n(b) calculating q=Bp=ΔδFp, where F is the circulant of kernel ƒ and Δδ is a flattening operator (whose form generally depends on the degree δ of the polynomial segments of F, and which operates on circulant F such that B=ΔδF is almost everywhere a locally constant matrix); and\n(c) backsolving Δδr=q to determine r=Fp. In one-dimensional cases (in which p has length N), Δδ has the form of the N×N circulant matrix:       Δ    δ    =      [                                        +                          (                                                                    δ                                                                                        0                                                              )                                                            -                          (                                                                    δ                                                                                        1                                                              )                                                            +                          (                                                                    δ                                                                                        2                                                              )                                                            -                          (                                                                    δ                                                                                        3                                                              )                                                ⋯                          0                                      0                                      +                          (                                                                    δ                                                                                        0                                                              )                                                            -                          (                                                                    δ                                                                                        1                                                              )                                                            +                          (                                                                    δ                                                                                        2                                                              )                                                ⋯                          0                                      ⋮                                                                                                                                                                                                                              ⋮                                                  -                          (                                                                    δ                                                                                        1                                                              )                                                            +                          (                                                                    δ                                                                                        2                                                              )                                                ⋯                                                                           0                                      +                          (                                                                    δ                                                                                        0                                                              )                                            ]  in which each entry is a binomial coefficient, and δ is the maximum degree of the spline segments of spline kernel ƒ. In higher-dimensional cases, the flattening operator Δ67  is defined similarly.\nIn other embodiments disclosed in U.S. application Ser. No. 09/480,908, the convolution Dp (where D is the circulant of smooth kernel d) includes the steps of:\n(a) approximating the kernel d by a spline kernel ƒ which is not a polynomial spline kernel (unless the kernel d is itself such a spline kernel, other than a polynomial spline kernel, in which case d=ƒ and step (a) is omitted);\n(b) calculating q=Bp=AFp, where F is the circulant of kernel ƒ and A is an annihilation or flattening operator, where A operates on circulant F in such a manner that AF=B is sparse when A is an annihilation operator, and A operates on circulant F in such a manner that AF=B is almost everywhere a locally constant matrix when A is a flattening operator; and\n(c) back-solving Ar=q to determine r=Fp.\nTo better appreciate the advantages of the present invention over conventional convolution methods, we next explain two types of conventional convolution methods: Fourier-based convolution and wavelet-based convolution.\nAs is well known, Fourier-based convolution relies on the elegant fact that if F is a Fourier matrix, sayFjk=e−2πijk/Nthen the transformation FDF−1 of the circulant is diagonal, whence we compute:Dp=F−1(FDF−1)Fp,where the far-right operation Fp is the usual Fourier transform, the operation by the parenthetical part is (by virtue of diagonality) dyadic multiplication, and the final operation F−1 is the inverse Fourier transform. For arbitrary D one requires actually three Fourier transforms, because the creation of the diagonal matrix FDF−1 requires one transform. However, if D is fixed, and transformed on a one-time basis, then subsequent convolutions Dp only require two transforms each, as is well known. The complexity then of Fourier-based cyclic convolution is thus O(N log N) operations (i.e., on the order of N log N multiplications and additions) for convolving a pattern p of length N (a pattern determined by N data values), because of the 2 or 3 FFTs (Fast Fourier Transforms) required. It should be noted that the Fourier method is an exact method (up to round-off errors depending on the FFT precision).\nAnother class of conventional convolution methods consists of wavelet convolution methods, which, by their nature, are generally inexact. The idea underlying such methods is elegant and runs as follows in the matrix-algebraic paradigm. Assume that, given an N-by-N circulant D, it is possible to find a matrix W (this is typically a compact wavelet transform) which has the properties:    (1) W is unitary (i.e. W−1 is the adjoint of W);    (2) W is sparse; and    (3) WDW−1 is sparse,where “sparse” in the present context denotes simply that any matrix-vector product Wx, for arbitrary x, involves reduced complexity O(N), rather than say O(N2).\nWith the assumed properties, we can calculate:Dp=W−1(WDW−1)Wpby way of three sparse-matrix-vector multiplications, noting that unitarity implies the sparseness of W−1. Therefore the wavelet-based convolution complexity is O(N) for convolving a pattern p determined by N data values, except that it is generally impossible to find, for given circulant D, a matrix W that gives both sparsity properties rigorously. Typically, if the convolution kernel d is sufficiently smooth, then a wavelet operator W (which is sparse) an be found such that within some acceptable approximation error the property (3) above holds. Above-noted properties (1) and (2) are common at least for the family of compact wavelets (it is property (3) that is usually approximate).\nAn advantage of “spline” convolution (in accordance with the teaching of U.S. application Ser. No. 09/480,908) over conventional wavelet convolution is that it can be performed (on data indicative of a pattern p comprising N data values) with cN arithmetic operations, whereas conventional wavelet convolution on the same data would require bN arithmetic operations, where (assuming typical error budgets) the factor “b” is significantly larger than the factor “c.” Among the other important advantages of the “spline” convolution method of application Ser. No. 09/480,908 (over conventional convolution methods) are the following: spline convolution is exact with respect to the spline kernel f, whereas wavelet convolution schemes are approximate by design (and error analysis for wavelet convolution is difficult to implement); the signal lengths for signals to be convolved by spline convolution are unrestricted (i.e., they need not be powers of two as in some conventional methods, and indeed they need not have any special form); and spline convolution allows acyclic convolution without padding with zeroes.\nSeparated-spline convolution in accordance with the present invention (like spline convolution in accordance with U.S. application Ser. No. 09/480,908) is an O(N) method for convolving a pattern p determined by N data values. Separated-spline convolution in accordance with the present invention has an advantage over spline convolution in accordance with U.S. application Ser. No. 09/480,908 in that separated-spline convolution in accordance with the invention can be performed (on data indicative of a two- or higher-dimensional patterns consisting of N data values) with dN arithmetic operations (multiplications and additions), whereas spline convolution in accordance with U.S. application Ser. No. 09/480,908 on the same data would require cN arithmetic operations, where the factor “c” is larger, and typically significantly larger, than the factor “d.” In other words, the implied big-O constant for separated-spline convolution according to the present invention is significantly smaller than the typical implied big-O constant for spline convolution as described in U.S. patent application Ser. No. 09/480,908."}
{"text": "In the building industry, rainwater gutters are used comprising a number of elongated channels connected end to end.\nEach channel extends in a given direction, has a bottom supporting face by which to mount the channel on a supporting surface, is bounded internally by a bottom wall defining a rainwater flow bed, and is bounded longitudinally by two substantially parallel end faces perpendicular to said direction.\nEach channel is normally connected to each adjacent channel by a number of bolts fitted through the corresponding end faces.\nKnown rainwater gutters of the type described above have several drawbacks, mainly due to the way in which the adjacent channels are connected not providing for any sort of scaling, and so allowing rainwater to seep between the end faces of the adjacent channels, thus damaging the construction beneath the gutter."}
{"text": "The pertinent prior art is shown and described in U.S. Pat. No. 3,806,280 having a coupling means comprising a U-shaped clevis that overlies the arm of the cradle and slides thereon to accomodate angular movement of the cradle to the linear movement of a control member."}
{"text": "The direct microscopy and the culture methods each have pitfalls In the past 20-25 years the direct visualization of bacteria in urine has largely been abandoned in favor of the methods involving culturing and counting the colonies of bacteria. Indeed virtually all of the studies of the significance of bacteriuria are based upon culturing the urine, and the direct microscopic examination of urine has been relegated to the status of a quick but inadequate screening procedure which may be helpful because it can be correlated with the culture methods.\nAny culture method requires that the bacteria will grow in the laboratory in the medium selected and in the time allotted. If the bacteria are damaged or dead when they left the body, then they will not grow There are many reasons why bacteria in urine would be damaged. The ionic strength or osmolality of the solution may be damaging. The oxidation potential of urine is usually too high (e.g., +0.22 to +0.25 volts). There may be a noxious metabolite in urine. (Human antibodies have been identified in urine and they have been demonstrated to be deposited on bacteria in urine.) Any or many of these factors may render a given bacterium non-viable in vitro. Finally, if the medium used is inappropriate for the growth of the particular organisms present, they will not grow. It can be readily shown by staining and microscopy that many of the bacterial forms found under the microscope were not alive at the time the specimen was obtained. For example, some do not contain nucleic acid, a biochemical component essential to life. Should all of the bacteria in a given specimen be devoid of nucleic acid, then none will grow and the culture of urine remains sterile. Indeed many urine specimens from sick patients containing huge numbers of bacteria will not yield a thriving bacterial culture in the hospital bacteriology laboratory. When the laboratory reports \"no growth\" the clinician may abandon the possibility of significant bacteriuria, and hence the possibility of an infectious cause. Nonetheless, these dead, damaged or fastidious bacteria, though they do not grow in culture, may in vivo have caused or exacerbated the patient's illness. FIG. 2 shows innumerable small cocci found, using the novel method taught herein, in the urine of a patient for whom the attempt to demonstrate bacteriuria by culture technique was completely unsuccessful. With regard to the direct examination of the urine, it must be noted that although bacteria may be seen in urine at only 100 diameters magnification, the size of the image is not the only consideration. Should the optical density and refractive index of a dead bacterium be near that of the medium, then it would not be detected by ordinary light microscopy. It may be seen by staining or perhaps by some specialized lighting. (Even then, as pointed out by Kunin, round bacteria cannot be distinguished from other near round particles such as crystals.) In my method the urine is examined wet at 100 to 400 diameters magnification, but it is also dried and prepared in a particular way so as to retain and preserve the bacterial structure through staining. In particular, I have found that urine contains lipids which act as detergents. Should they be allowed to remain on the slide when an aqueous dye is applied to the slide, then much of the sediment (including bacteria) will wash off of the slide and the preparation will be lost. This is a major reason why past attempts to study bacteriuria have failed. The photograph of FIG. 1 illustrates the difference. In that photograph one slide had been prepared in the standard way and the other has been washed with a lipid solvent. After staining much more sediment is found on the washed slide. Most of the sediment had washed off in the standard preparation. A chromatogram of the lipids removed reveals several lipids in the range of polarity of the phospholipids (e.g., lecithin, phosphatidylserine, etc. which substances act as detergents) but they do not contain appreciable phosphorus and thus they are not phospholipids. Standard methods of preparing and staining urine specimens, such as that of Melnick, U.S. Pat. No. 4,225,669, do not provide for the precautionary removal of these lipids.\nRheumatoid arthritis (RA) is a chronic inflammation of the joints, generally regarded as a systemic autoimmune disorder. Its etiology is unknown, but it has been postulated that it is associated with microbial infection. See, e.g., D. C. Demonde, ed., Infection and Immunity in the Rheumatic Diseases, 95-287 (Blackwell Scientific Publications, London: 1976). The evidence, however, until the present discovery, was inconclusive. See, e.g., D. J. McCarty, et al., ed., Arthritis and Allied Conditions: A Textbook of Rheumatology, ch 28 at 417 (9th ed. 1979); R. G. Petersdorf, et al., ed., Harrison's Principles of Internal Medicine, Part Six, Chapter 346, at 1977 (McGraw Hill: 1983). Bacteriuria has not been associated with RA, and indeed one authority remarks \"Urinary abnormalities are relatively uncommon in RA . . . Urinary tract infection was not found to be increased in RA patients.\" McCarty, supra, chapter 33, page 499, citing Ann. Rheum. Dis., 27: 345 (1968). Hypertension is a chronic elevation of blood pressure resulting from the obstruction of blood flow within the kidney (secondary hypertension) or without apparent cause (essential hypertension). One kidney disorder associated with secondary hypertension is pyelonephritis, the inflammation of the renal pelvis of the kidney as a result of bacterial infection, usually responsive to antibiotics. It has not been reported, however, that there is any correlation between essential hypertension and asymptomatic bacteriuria (bacteriuria observed in patients not reporting symptoms of urinary tract disorders). According to N. M. Kaplan, Clinical Hypertension, 14 (3d. ed. 1982), bacteriuria is found in 2-5% of hypertensives. Most of these positive cultures were of gram-negative rods. The method of the present invention has demonstrated a much higher incidence of bacteriuria in hypertensives, perhaps as high as 90%, and that cocci or \"exploded cocci\" are found in considerable numbers."}
{"text": "Thin parts are, for example, waterproof rubber or the like to be waterproof stoppers of take-in openings for electric wires which are attached to terminals of electric wires and are provided to terminal housings. Mounting apparatuses that insert terminals of electric wires separately prepared by cutting out into the waterproof rubber so as to mount the waterproof rubber are prepared for such thin parts (waterproof rubber). The electric wires and the waterproof rubber as the thin parts are sequentially supplied to the mounting apparatuses, and a lot of electric wires which are inserted into the waterproof rubber are produced and are used for wiring works.\nWaterproof rubber X as the thin part has various shapes according to thickness of electric wires and intended use. For example as shown in FIG. 14, the waterproof rubber X has an approximately cylindrical shape, and a hole H through which electric wires are inserted is pierced in the waterproof rubber X. A longitudinal center of the cylinder has a step portion P1, and both ends P2 and P3 have different thickness. A feeding device that takes out pieces of such waterproof rubber one by one so as to feed them is provided for the waterproof rubber, and thus, concretely the waterproof rubber is sequentially fed to the mounting apparatus for the electric wires, so that the electric wires are inserted into the waterproof rubber. Such thin parts include various ones having length of 7 mm and thickness of 5 mm and having length of 7 mm and thickness of 1.2 mm.\nThe conventional waterproof rubber feeding device is, for example as shown in FIG. 15, provided with a receiving portion 502 in a position where the waterproof rubber X, which is put into a rotating drum 500, is lifted by a blade portion 501 on a drum inner wall and then is dropped, can be received. A transporting device which is called as a feeder 503 which follows the receiving portion transports the waterproof rubber to a next step. The feeder 503 is the transporting device composed of a groove-shaped rail, and the waterproof rubber received by the groove of the rail is advanced by vibration so as to be fed.\nIn the above feeding device, excessive pieces of waterproof rubber are accumulated on the receiving portion, and thus the waterproof rubber is occasionally prevented from being fed to a next step. For this reason, there is Patent Document 1 (Japanese Patent Application Laid-Open No. 2003-20115) by the applicant of the present invention which is filed in order to solve this problem. According to this, since a quantity of the waterproof rubber which is fed to the receiving portion can be adjusted suitably, the effect that the receiving portion is not blocked up is obtained.\nFurther, in Patent Document 2 (Japanese Utility Model Laid-Open No. 02-8198), a chip part (thin part) which is fed from an air hopper drops onto a rail and enters a groove on the rail according to the part size. Excessively fed thin parts bump against the stopper and are collected to a parts pocket. The thin parts which enter the groove of the rail are adapted to be fed by applying slight vibration to the rail by means of a vibrator fixed under the rail.\nIn such a manner, the thin parts are received by the receiving portion and are transported to the feeder, but at this time the thin parts having the approximately cylindrical shape are stored in a groove of the feeder in an upright posture, and are transported sequentially. At this time, however, when the thin parts having top-bottom directional properties are fed to a next step, their vertical directions should be aligned. It is, therefore, necessary to remove thin parts whose direction is opposite to a desired direction from the feeder by means of a certain method and feed only thin parts which face the desired direction.\nAs mentioned above, in conventional methods, devices tend to be large because a space for dropping the thin parts onto the receiving portion is required, a space for again collecting thin parts which drop without being received by the receiving portion is required, or the thin parts should be lifted to a place where they are dropped. Further, when oil is applied to the thin parts (waterproof rubber), the feeding device is comparatively easily contaminated, and thus it should be suitably cleaned. It is, however, difficult to clean the conventional feeding devices due to their constitution. Further, when a difference in dimension between ends of the thin parts is small or discrimination between top and bottom portions is difficult, disadvantageously the thin parts are not fed due to their shapes. The thin parts (waterproof rubber) fed from the hopper to the feeder whose direction is opposite to a desired direction should be removed, and this is inefficient."}
{"text": "In order to maintain or grow their market share, manufacturers of disposable absorbent articles such as disposable diapers and absorbent pants must continue to discover and develop improvements to materials, components and features that affect aspects such as containment, absorbency, comfort, fit and appearance. Absorbent pants are manufactured in smaller sizes to be used as, e.g., pull-on diapers and toilet training pants for young children, and in larger sizes to be used as, e.g., undergarments for persons suffering from incontinence.\nA particular type of absorbent pant design currently marketed is sometimes called the “balloon” pant. The balloon pant design usually includes a central absorbent chassis and an elastic belt. The elastic belt is usually relatively wide (in the longitudinal direction) and elastically stretchable in the lateral direction. It entirely encircles the wearer's waist, and thereby covers a relatively large amount of the wearer's skin, and also makes up a relatively large portion of the visible outside surfaces of the pant. The central chassis portion is typically joined to the inside of the belt in the front, wraps under the wearer's lower torso between the legs, and is joined to the inside of the belt in the rear. The belt is often formed of two layers of nonwoven web sandwiching one or more elastic members such as a plurality of laterally-oriented strands or strips of elastomeric material, or a section of elastomeric film or elastomeric nonwoven. It is common among such designs that, in manufacture, the elastic member(s) are sandwiched between the nonwoven web layers in a strained condition. Upon completion of manufacture and allowance of the belt with sandwiched elastic member(s) to assume a relaxed condition, the elastic member(s) contract laterally toward their unstrained lengths. This causes the nonwoven web layers to form gathers that take the form of visible rugosities in the belt. It is believed that some consumers find the rugosities attractive because they present a textured, plush, frilly and/or soft appearance and feel, and also provide a visible indication of stretchability and comfortableness.\nThe balloon pant design provides for certain efficiencies in manufacture, and it is believed that the design may gain popularity. Consequently, any improvement in components such as the belt that enhance its appearance and functionality may give the manufacturer thereof an advantage."}
{"text": "In general, an automatic door opening and shutting device detects either person(s) or vehicle(s) coming in and out the door by a detecting means using a limit switch or a photoelectric tube. And by actuating an opening and shutting means with detecting signals, the door is opened or shut automatically. In case the automatic door is opened or shut by a friction drive system, the disadvantage is that power transmission on the friction surface tends to become irregular due to some shakes of the door being caused during the door opening or shutting, and then the automatic door is not actuated smoothly.\nFurther, in case the automatic door is switched to a manual operation by some accidents and other causes, it is forced to be opened or shut with a stronger force because a load of a friction drive means is applied to the friction surface of the door. Accordingly, it is not easy to open or shut smoothly by hand the door which ceased its automatic function.\nFor the sake of removing the aforegoing defects and disadvantages, the present invention has been achieved."}
{"text": "The somatolactogenic hormones prolactin (PRL) and growth hormone (GH) are necessary for the full growth and maturation of vertebrate species.\nProlactin (PRL) was originally identified as a neuroendocrine hormone of pituitary origin. PRL expression has also been detected in the decidua, breast and T-lymphocytes (Clevenger, C. V. and Plank, T. L. J. Mammary Gland Biol. Neoplasia 1997 2:59–68; Mershon et al. Endocrinology 1995 136:3619–3623; DiMattia et al. Endocrinology 1986 122:2508–2517; Ginsburg, E. and Vonderhaar, B. K. Cancer Res. 1995 55:2591–2595; Gellersen et al. Mol. Endocrinol. 1994 8:356–373; Clevenger et al. Proc. Natl. Acad. Sci. USA 1990 87:6460–6464; Montogomery et al. Biochem. Biophys. Res. Commun. 1987 145:692–698). A primary function of this hormone lies within the breast. However, functional pleiotropism of this peptide with regard to reproduction, osmoregulation and behavior has also been recognized (Nicoll, C. S. Handbook of Physiology; Section 7: Endocrinology, pp. 253–292, Washington, D.C.: American Physiology Society. 1974). Several lines of evidence have also indicated an immunoregulatory role for this peptide (Clevenger et al. Journal of Endocrinology 1998 157:187–197; Weigent, D. A. Pharmacol. Ther. 1996 69:237–257). Structural analysis of PRL has revealed it to be related to members of the cytokine/hematopoietin family which also includes growth hormone, erythropoietin, granulocyte-macrophage colony stimulating factor (GM-CSF) and interleukins 2–7 (Bazan, J. F. Immunol. Today 1990 11:350–354).\nThe pleiotropic actions of PRL are mediated through its receptor (PRLr), a member of the superfamily of type I cytokine receptors. PRLR is present on numerous tissues including mammary epithelia, T and B lymphocytes and macrophages (Dardenne et al. Endocrinology 1994 134:2108–2114; Pellegrini et al. Mol. Endocrinol. 1992 6:1023–1031). Acting through its receptor, PRL signaling stimulates cell proliferation, survival and cellular differentiation in a tissue- and microenvironment-dependent manner. In the mammary and immune systems, PRL is believed to act at the endocrine, paracrine, and autocrine levels in regulating T-lymphocyte proliferation and survival (Gala, R. R. PSEBM 1991 198:513–527; Yu-Lee, L. Y. Proceedings of the Society for Experimental Biology and Medicine 1997 215:35–52; Kooijman et al. Adv. Immunol. 1996 63:377–454; Prystowski, M. B. and Clevenger, C. V. Immunomethods 1994 5:49–55) and the terminal maturation of mammary tissues (Kelly et al. Rec. Prog. Horm. Res. 1993 48:123–164; Shiu et al. Rec. Prog. Horm. res. 1987 43:277–289). PRL is also believed to act as both an endocrine and autocrine/paracrine progression factor for mammary carcinoma in both rodents and humans (Welsch, C. W. Cancer Res. 1985 45:3415–3443; Welsch, C. W. and Nagasawa, H. Cancer Res. 1977 37:951–963; Manni et al. Cancer Res. 1986 37:951–963; Malarkey et al. J. Clin. Endocrinol. Metab. 1983 56:673–677; Clevenger et al. Am. J. Pathol. 1995 146:695–705; Fields et al. Lab. Invest. 1993 68:354–360; Ormandy et al. J. Clin. Endocrinol. Metab. 1997 82:3692–3699; and Mertani et al. Int. J. Cancer 1998 79:202 22 79:202–211).\nPleiotropic actions of GH are also largely mediated through a type I cytokine receptor, GHr.\nLigand-induced dimerization of PRLr and GHr activates several associated signaling cascades including the Jak-Stat, Ras-Raf, and Fyn-Vav pathways (Campbell et al. Proc. Natl. Acad. Sci. USA 1994 91:5232; Clevenger et al. J. Biol. Chem. 1994 269:5559; Clevenger et al. Mol. Endocrinol. 1994 8:674; Clevenger et al. J. Biol. Chem. 1995 270:13246). However, studies indicate that both PRL and GH are internalized via an endosomal-like pathway and transported across the endoplasmic reticulum (ER) and nuclear envelopes (Clevenger et al. Endocrinology 1990 127:3151; Rao et al. J. Cell Physiol. 1995 163:266). This process is referred to as nuclear retrotranslocation. The mechanism of this retrotranslocation, and the nuclear action of these somatolactogenic hormones, however, is not well understood.\nBoth PRL and GH lack enzymatic activity. These hormones also contain no nuclear translocation signal. Thus, for PRL and GH to act within the nucleus, they must do so through a binding partner or chaperone.\nCypB is a member of the cyclophilin family of cis-trans peptidyl prolyl isomerases (PPI) (Price et al. Proc. Natl. Acad. Sci. USA 1991 88:1903; Ruhlmann, A. and Nordheim, A. Immunobiol. 1997 198:192; Resch, K. and Szamel, M. Int. J. Immunopharmac. 1997 19:579). This family of proteins was initially identified as the binding partners for the immunosuppressive agent cyclosporine (CsA). CsA interacts with the cyclophilin with high affinity, inhibiting their PPI activity and the action of the phosphatase calcineurin, necessary for NF/AT-transactivated expression of IL-2 (Kronke et al. Proc. Natl. Acad. Sci. USA 1984 81:5214; Liu et al. Cell 1991 55:807 66:799; Friedman, J. and Weissman, I. Cell 1991 66:799; McCaffrey et al. J. Biol. Chem. 1993 268:3747; Bram, R. J. and Crabtree, G. R. Nature 1994 371:355; Bram et al. Mol. Cell Biol. 1993 13:4760). Structurally CypB is a β-barrel protein containing both N-terminal ER-leader and putative nuclear translocation signal sequences and C-terminal ER-retention sequences (Allain et al. J. Immunol. Meth. 1995 178:113; Mariller et al. Biochem. Biophys. Acta 1996 1293:31). CypB has been observed in the ER and nucleus, and can be found in appreciable levels in blood (150 ng/ml) and breast milk (Hirada et al. Cell 1990 63:303; Price et al. Proc. Natl. Acad. Sci. USA 1994 91:3931). Cyclophilins, via their PPI activity, facilitate protein folding and have been shown to contribute to the maturation of several proteins, including carbonic anhydrase and the HIV glycoprotein Gag (Taylor et al. Prog. Biophys. Molec. Biol. 1997 67:155; Streblow et al. Virology 1998 245:197). Despite these insights, the physiologic function of CypB has remained uncertain.\nIt has now been found that cyclophilin B (CypB) interacts specifically with somatolactogenic hormones, PRL and GH, as a chaperone mediating the transport, maturation and/or function of these proteins."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a flash memory and a method of forming a flash memory, and more specifically a method of forming a flash memory cell using an asymmetric control gate with a sidewall floating gate.\n2. Description of the Related Art\nA flash memory is unique in providing fast compact storage which is both nonvolatile and rewritable.\nIn a flash memory, the threshold voltage Vt for conduction of a field effect transistor (FET) changes state depending upon the amount of charge stored in a floating gate (FG) part of the FET. The floating gate is a charge storing region which is isolated from a more traditional gate conductor CG (control gate or \"wordline\") by a thin dielectric. The states of the Vt change with the amount of charge stored by the FG.\nSince the FG directly controls conductivity between source and drain in a channel, the state of a FG memory cell is determined by applying certain voltages to the source or drain of the FET and observing whether the FET conducts any current.\nFlash memory cells with a sidewall floating gate occupy a smaller area than those with conventional (layered) floating gates. For example, in U.S. Pat. No. 5,115,288, sidewall gates were formed on one side of the wordline by employing an extra mask. Sidewall spacers were formed on both edges of the wordline, then removed along one of the edges using the extra mask and an etching operation. Thus, the conventional approach uses a trim mask to define the floating gate.\nHowever, this approach is expensive and requires good control of the overlay for the spacer removal mask.\nOther conventional structures also are known. For example, in one structure, polysilicon spacers formed on both sidewalls are used for the floating gate. One spacer sits on top of the tunnel oxide area for programming. The other spacer is called an \"added-on floating gate\". Both spacers are linked by a polysilicon body. However, a large cell size results.\nIn a second conventional structure, only one polysilicon spacer is used as the floating gate. A mask must be aligned to the top of the control gate, to remove the other floating gate spacer. Hence, the control gate cannot be small, since, otherwise, any misalignment will cause a problem. Therefore, this cell has difficulty in being down-scaled.\nIn yet another conventional structure, similar to the second conventional structure described above, a mask is needed to remove a sidewall spacer floating gate. Further, this spacer has a re-entrant corner which is very difficult to be completely removed.\nThus, the conventional methods require extra process steps, material and more precise lithographic alignment, thereby resulting in increased manufacturing costs."}
{"text": "The present invention relates generally to a circuit arrangement used as an interface for a sensor, in particular, for a pumped reference oxygen sensor used in combination with a combustion engine.\nOxygen sensors are particularly used in combination with combustion engines using controlled catalytic converters. The use of catalytic converters in cars started much earlier in the USA in comparison with Europe. However, innovation cycles are often slower in the US automotive technology. Therefore, often older technology is used for longer periods in the USA as compared to European countries. For example, the binary lambda oxygen sensor which is used to regulate the gasoline mixture in a combustion engine, comprises the negative terminal being electrically coupled with the sensor housing in embodiments of the first generation. The reasons for this connection relates to an easier and cheaper construction. However, this construction also results in an electrical coupling between the sensor housing and the engine through the fixture of the sensor within the muffler arrangement. This connection is disadvantageous because shifted ground potentials within a motor vehicle. For example, the engine ground is usually more negative than the ground of the motor control unit. This potential shift is due to switching of high load currents, e. g., 10 . . . 15A, and the inner resistance of the ground wires, e. g. 20 . . . 40 mΩ. Typical voltage shifts are in the range of −300 . . . +600 mV. In addition, any on/off switching of high loads causes voltage overlay peaks of up to multiple 100 m Vss. These ground distortions and overlay voltages can cause serious problems with respect to evaluation of the respective sensor signals and may render the actual sensor signals completely useless.\nOxygen sensors according to newer technology are, therefore, fully isolated and, thus, avoid any electrical coupling with the engine. However, oxygen sensors of the first generation are still widely used, in particular, in the United States because of their lower manufacturing costs. Therefore, modem motor control units must be able to interface with these kind of sensors which are not fully isolated. To this end, specific interface circuits used to be available which allowed for the evaluation of sensors whose housing is electrically coupled with the engine and thus with the motor vehicle ground.\nFIG. 4 shows an example of such a known circuit. A sensor 400 which is connected with the engine ground is coupled through resistors 430 and 415 with the input terminals of the interface circuit proper. Capacitors 425 and 435 are coupled between the input terminals of the interface circuit and the interface ground. One of the sensor connections is furthermore coupled through resistor 410 with a supply voltage Vcc and through capacitor 420 with the interface ground. The interface circuit comprises a first switch 440 which is coupled with the input terminals. The switch output is coupled through a capacitor 445 with the input of a second switch 460 and through capacitor 450 with the input of a third switch 455. The first output of switch 455 is connected with the supply voltage Vcc. The second output of switch 455 and the first output of switch 460 are coupled with the interface ground. The second output of switch 460 is connected to the inverting input of an integrator consisting of operational amplifier 475 and capacitor 470 in its feedback loop. A fourth switch couples the inverting input of the integrator with either the supply voltage or the output of the integrator. The output of the integrator is coupled with the first input of a fifth switch 495 which is controlled by the output signal of a comparator 405. The second input of switch 495 is coupled through resistor 490 with the supply voltage. The first input of comparator 405 is coupled with the second interface input terminal and the second input of comparator 405 receives a voltage signal being equal to half the supply voltage. Furthermore, a timing circuit is provided which generates control signals for switches 440, 455, 460, and 480.\nAll switches are implemented as CMOS switches. Capacitor 445 is used as a transfer capacitor for eliminating the common mode of the input signal. To this end, the capacitor is switched in a first position between interface ground and the second input terminal and in a second position between the inverting input of operational amplifier 475 and the first interface input terminal with a high frequency. Thus, the CMOS switch operates like a resistor. Capacitor 465 operates as a feedback capacitor in a similar way. These two capacitors operating as resistors form together with the operational amplifier/integrator an inverting amplifier. The bias capacitor 450 together with CMOS switch 455 are used to generate a small bias current which is fed to the sensor 400 and which will not influence the measurement when the sensor is in operating mode, i.e. the sensor has low resistance.\nOne of the disadvantages of this circuit arrangement is that the CMOS switches at the input of the circuit must comply with a high standard. This renders this circuit expensive and interference-prone. In addition, this circuit must withstand the required negative input voltages. thus, additional protective measurements, such as, isolation and charge pumps (not shown) must be provided. Furthermore, the CMOS switches must be able to tolerate a relatively high input voltage of up to 12V in case of a short circuit of the sensor. This is particularly difficult because the supply voltage is usually only 5V. Integrated circuits using this technology need furthermore additional isolation/separation measurements if more than one interface circuit is provided to prevent any cross over influence of the channels and to prevent a latch-up.\nThe bias current generated by switch 455 and capacitor 450 is used to detect a connection failure between the first input terminal and the sensor. In such a case, the bias current will overdrive the operational amplifier. A similar scenario takes place in case of a short circuit between the first input terminal and the positive terminal of the battery. The output voltage in both cases will be approximately 0V. To detect any interruption between the second input terminal and the sensor additional circuitry is necessary. This additional circuitry is shown in FIG. 4 with resistor 410, capacitor 420, comparator 405 and CMOS switch 495. During normal operation the current generated by resistor 410 will flow to the engine ground through the electrical coupling of the housing of sensor 400 and will not influence the measurement. However, in case of an interruption of this connection the potential at the second input terminal will raise to the supply voltage, for example, 5V. In case of a normal operational temperature of the sensor, i.e. low resistance of the sensor, the operational amplifier will be driven to its positive limit, e.g., 5V. However, in case of a cold sensor (during the start up phase of the engine) the sensor will have a high resistance and the internally generated bias current will put the circuit into an undefined state. To prevent such a state, comparator 405 will compare the potential at the second input terminal with Vcc/2. If the potential is above this threshold, comparator 405 will control switch 490 to select a constant output voltage to signalize this error.\nAs described above, the prior art interface circuit is highly cumbersome and requires additional evaluation of the generated output signal. Furthermore, this type of interface circuit is not in production anymore and, thus, not available for new construction which specifies the use of a non-isolated sensor."}
{"text": "Generally, many techniques for signal compression and recovery have been introduced so far. And, applicable targets of the corresponding techniques are various data including audio data, video data and the like. Moreover, the signal compression or recovery techniques evolve in a direction to enhancing audio or video quality with high compression rates. Besides, many efforts have been made to raise transmission efficiency for adaptation to various communication environments.\nHowever, it is still believed that there exists a margin for the enhancement of the transmission efficiency. So, many efforts need to be made to maximize transmission efficiency of signals in the very complicated communication environments through the development of new processing schemes for signals."}
{"text": "This invention relates to a connector which is provided in an external device, such as a card reader for a game-program IC card, and which is adapted for connection to the IC card.\nA conventional connector provided in an external device for an IC card is constructed as shown in FIGS. 1A, 1B, and 1C. An 1C card which is inserted into this connector has a structure such as that shown in FIGS. 2A, 2B, and 2C. As shown in these figures, a card reader connector 1 which is a connector disposed in an external device (not shown) for the card has a card-insertion opening 2 and a card-receiving portion 3. An IC card 4 is formed such that it can be inserted into and pulled out of the card-receiving portion of the card reader connector. The IC card 4 has a package 5 in which a semiconductor device 4a is encapsulated, metal panels 5a disposed on upper and lower surfaces of this package, and a plurality of external terminals 6 arranged on one end of the package 5. The package 5 is formed of an insulating material, and an insulating material coats all surfaces of the metal panels 5a. A card shutter 7 is provided on the package 5 of the IC card 4 so as to be slidable in an insertion direction of the IC card and adapted to protect the external terminals 6 when the card is not used. The card reader connector 1 is provided with shutter-opening-closing pins 8 adapted to open the shutter 7 by moving it in the direction opposite to that of the insertion of the card, and a plurality of electrodes 9 capable of being elastically deformed to contact with the external terminals 6 respectively when the card is inserted into the connector. The shutter-opening-closing pins 8 can be made of metal or an insulating material. The shutter 7 is provided with slide members 10 having pin-contact portions 10a (see FIG. 2B) which contact the pins 8 when the IC card 4 is inserted into the card reader connector 1.\nWhen the IC card 4 is used with the thus-constructed card reader connector 1, the IC card 4 is inserted through the card-insertion opening 2 into the card-receiving portion 3 of the card reader connector 1. At the same time, as shown in FIGS. 3A, 3B, and 3C, the pin-contact portions 10a of the slide member 10 are brought into contact with the pins 8 so as to open the shutter 7 and to connect the external terminals 6 of the IC card 4 and the electrodes 9 of the card reader connector 1 to each other.\nThe connection of the electrodes 9 of the card reader connector 1 and the external terminals 6 of the IC card 4 is released by pulling the IC card 4 out of the card-receiving portion 3.\nIn this type of connector for connecting an external device to the IC card, friction between the card reader connector 1 and the IC card 4 and/or between the shutter 7 and the package 5 of the IC card 4 tends to generate static electricity when the IC card 4 is inserted into or pulled out from the card reader connector 1. If this static electricity flows through the external terminals 6, there is a possibility of the semiconductor device 4a in the IC card 4 being damaged or a possibility of software errors occurring due to noise, resulting in reduction of the reliability of the IC card operation."}
{"text": "A head mounted computer has a display and is mounted to a user's head to enable the user to view the display. The display shows images to the user, for example, from a computer or a remote device. The user can control the head mounted computer or another remote device, which in turn affects the images shown on the display."}
{"text": "Recently, studies on technologies of efficiently capturing, transporting and storing carbon dioxide from off-gas streams in order to reduce greenhouse gas emissions have been actively conducted. At present, absorption processes that use wet absorbents are being widely used to separate and remove carbon dioxide from off-gas streams. In general absorption processes, absorbents such as amine-based compounds are used to absorb and separate carbon dioxide from off-gas streams. However, this absorption process has a disadvantage in that the amine-based compounds are expensive, and thus the costs of designing and operating an apparatus for separating carbon dioxide increase. In addition, the absorption process has problems in that the consumption of energy for regeneration of the absorbent is high and an apparatus for treating off-gas is corroded by the amine-based compound."}
{"text": "Wearable computing devices are electronic devices worn by a user. These devices may be worn differently by the user based on the device's intended function. For example, wearable computing devices to track a user's location or activity level may be worn similar to a bracelet, wristwatch or any other similar appropriate form factor. Heads up display (HUD) devices may have a form factor similar to a pair of eyeglasses.\nOne of the disadvantages of current wearable computing device design is that necessary power and input/output (I/O) interconnections are easily visible, and thus the device cannot completely have the equivalent aesthetic and form factor of a traditionally worn accessory. Furthermore, accessories such as bracelets and watches require some type of fastener (i.e., a clasp) for the user to put on or take off the accessory. This “break” in the wearable device often prevents interconnectivity and power from being provided throughout the device.\nDescriptions of certain details and implementations follow, including a description of the figures, which may depict some or all of the embodiments described below, as well as a discussion of other potential embodiments or implementations of the inventive concepts presented herein. An overview of embodiments of the invention is provided below, followed by a more detailed description with reference to the drawings."}
{"text": "This invention relates to surgical tool holders, and more particularly, to holders of rotary surgical cutting tools.\nSurgical tools and their respective holders have to be kept clean and sterile before any use in a hospital environment in order to minimize risk of transfer of disease or infection from patient to patient following the emergence of certain “prions” that are not killed by normal hospital sterilization and thus need to be physically removed through washing and rinsing. A surgical tool, for example for preparing for the fitting of a hip prosthesis, works in a medium which causes considerable soiling of the tool and of the tool holder. Despite the importance of doing so, the thorough cleaning of these devices is difficult. Surgical tool holders of the prior art are designed such that, in certain cases, washing and rinsing is generally not an effective way of cleaning the tool as some of the more resilient bugs are not destroyed. This is due to the small spaces left between component parts which allow only minimal access by cleaning agents.\nFurther, small diameter surgical cutting tools of the prior art are generally of solid construction. When cutting, such tools do not provide for the collection of bone chips for use in grafting.\nWhat is needed therefore is a surgical tool holder which is quickly and simply disassembled for cleaning and sterilization. Still further, what is needed is a small diameter hollow tool that provides for the collection of bone chips for grafting."}
{"text": "The present invention generally relates to athletic garments and, more particularly, to an athletic garment with an integral cup or cup assembly for holding an athletic protective device and/or an incontinent pad.\nIn many athletic events, players may suffer injuries in their groin regions as a consequence of impacts to that region. To avoid, or at least minimize injuries in the male groin region, players use different types of protective devices, such as athletic protective cups. Athletic protective cups are typically used in many contact sports, such as football, baseball, hockey, soccer and the like. Protective cups, however, are also used in non-athletic events where the users are susceptible to impacts to the groin region.\nProtective cups should be secured against the user's body to protect the groin region from impact. Several apparatus and devices have been developed over the years to secure a protective cup over a male's genitals. For instance, athletes sometimes use jockstraps to hold a protective cup in place. U.S. Pat. No. 4,134,400 describes an example of a jockstrap. This jockstrap includes a waistband, multiple straps attached to the waistband, and a pouch for holding the protective cup. The pouch is attached to the waistband and to all the straps. The straps are specifically attached to the back portion of the waistband. While the type of jockstrap disclosed in U.S. Pat. No. 4,134,400 can secure a protective cup over the user's groin region, it can be confusing and difficult to put on due to the multiple straps. In order to wear the jockstrap, the user must place his legs through specific straps. However, the user may easily put his legs through the incorrect straps. If the jockstrap is not worn properly, the protective cup may move out of its proper position, leaving the user unprotected from blows to the groin. In addition, the straps of the jockstrap may stretch due to regular wear and tear, becoming incapable of holding the protective cup in the proper position.\nInstead of a jockstrap, athletes may use a jock short or brief to secure a protective cup against their bodies. U.S. Pat. Nos. 5,157,793 and 5,819,323 describe jock shorts or briefs capable of holding a protective cup and made of a stretchable fabric, such as the polyurethane-polyurea copolymer fabric, sold under the trademark LYCRA® by Invista North America S.A.R.L. Corporation. These kinds of shorts tightly fit the user's body. Nevertheless, these shorts do not provide enough strength to hold and support the protective cup over the user's groin region. In this type of shorts, gravity and the constant movement by the user may cause the cup to move out of its proper place. In such case, the user may be injured if he receives a blow to his genitals.\nOther kinds of jock support shorts may fit loosely the user's body and rely on straps to hold a protective cup in place. U.S. Pat. Nos. 7,757,307 and 7,216,371 describe shorts, which employ straps to support a protective cup. In these types of shorts, the straps are attached to a pocket, which holds the protective cup. These straps suffer from the same shortcomings as the straps of the jockstrap described above. Thus, during use, the protective cup may move out of its proper position.\nIn light of the shortcomings of current devices for securing a protective cup in its proper place, there is a need for a device, apparatus or garment capable of properly securing and holding a protective device over a wearer's genitals."}
{"text": "This invention relates to the domain of advanced driver assistance systems and autonomous robots. Driver assistance systems are automated technical devices, typically involving sensors such as cameras, laser scanners or infrared sensors, which are mounted to vehicles (cars, motorbikes, planes, boats, trains, . . . ). Their purpose is to monitor the environment of the vehicle thus equipped and assist the driver while driving in real-world traffic environments. “Assisting” thereby can be the control of an actuator of the vehicle (brakes, steering, accelerator, . . . ) or at least the generation of visual or audio signals assisting the driver.\nExamples include lane keeping systems, braking assistants and distance monitoring devices, just to name a few. The invention proposes a driver assistance system that focuses on visual processing and is therefore equipped e.g. with a stereo camera as sensor device."}
{"text": "The invention relates to manufacturing an inner conductor of a resonator.\nResonator structures of a high frequency area, a radio frequency area in particular, are used e.g. in base stations of mobile telephone networks. Filters may utilize resonator structures e.g. as adapting and filtering circuits in transmitter and receiver units of the base stations.\nA resonator structure comprises an inner conductor of the resonator attached to an attachment surface, which in practice most often is an end, such as a bottom or a cover, of a housing structure serving as an outer conductor of the resonator structure. The inner conductor is thus short-circuited to the attachment surface, i.e. in practice to the outer conductor. A short-circuited end of the inner conductor, at which the inner conductor is thus short-circuited to the outer conductor, is also called an inductive end owing to the fact that signal coupling at the short-circuited end is mainly carried out inductively.\nAt a second end of the inner conductor, the inner conductor is galvanically separated from the outer conductor, so this end is the xe2x80x9cfreexe2x80x9d end of the inner conductor. The free end of the inner conductor is also called a capacitive end of the inner conductor owing to the fact that signal coupling at this end is mainly carried out capacitively. The outer conductor and the inner conductor located within a section defined by the outer conductor together form a resonance circuit. In practice, the resonator structures often comprise a plurality of circuits, i.e. the resonator structure comprises several pairs comprising an inner conductor and an outer conductor, i.e. each section formed by the outer conductor comprises a separate inner conductor. The resonance circuits of a multi-circuit resonator structure together form a desired frequency response for the resonator structure.\nNormally in a coaxial resonator, the inner conductor of the resonator is a straight wire or a pin attached only to the bottom of the resonator. Such a resonator is long and thus takes a lot of space. The resonator pin is quite easy to manufacture. The problem then is, however, how to adjust the coupling of the resonator since it is difficult to attach such a controlling element to the resonator pin that would enable the resonator to be easily coupled to e.g. an adjacent resonator. Furthermore, the capacitive coupling provided by the wire-like inner conductor is poor.\nIn order to decrease the space required by the resonator, for instance a helix coil is used as the inner conductor, in which helix coil the same operational length fits into a shorter space since the resonator in the helix resonator is formed as a coil. The helix coil is, however, difficult to manufacture. A further drawback is that it is extremely difficult to attach to the helix coil a coupling wire or other such projection necessary when the coupling between two resonance circuits is to be adjusted. A further problem with the helix resonators is the difficulty to support them and carry out the temperature compensation. An inner conductor implemented by utilizing a helix coil cannot provide a high-quality capacitive coupling.\nA known solution for controlling the resonance frequency of a resonator circuit is a solution wherein an adjuster bolt located in the cover of a filter serves as the frequency controlling element, and the distance of the adjuster bolt with respect to the free end of the resonator located in a section under the cover is adjusted by turning the bolt. The solution is not the best possible one since it requires additional structures on the outer surface of the housing. A further problem is that the adjuster bolt requires that the cover of the filter should be thick or the cover should at least comprise a thicker section to enable threads to be provided on the cover for the adjuster bolt, or, alternatively, to enable a nut-like part with threads attached to the cover to be used. The cover has to be thick particularly because it also needs to be rigid in order to prevent the distance of the frequency controlling element in the cover with respect to the resonator from changing after the controlling procedure and from further causing the capacitance, and thus the resonance frequency, to change in an undesired manner.\nAn object of the invention is thus to provide a method of manufacturing an inner conductor of a resonator, and an inner conductor so as to enable the above-mentioned problems to be alleviated. This is achieved by a method disclosed in the introduction, characterized by manufacturing at least part of the inner conductor from a uniform, electrically conductive material blank by utilizing a deep-drawing method wherein the blank is struck or pressed with a tip of an impact device, whereby during each stroke or pressing, the tip draws more and more blank material in the direction of the stroke.\nThe invention further relates to an inner conductor of a resonator comprising a first end and a second end, which is free.\nThe inner conductor of the invention is characterized in that at least part of the inner conductor is deep-drawn from a uniform, electrically conductive blank.\nPreferred embodiments of the invention are disclosed in the dependent claims.\nThe idea underlying the invention is that the inner conductor is manufactured by utilizing a deep-drawing method.\nSeveral advantages are achieved by the method and inner conductor of the invention. The deep-drawing method enables the inner conductor and a flange located at the free end thereof to be manufactured virtually simultaneously. In addition, a potential projection or a site for the same can be manufactured in connection with manufacturing the inner conductor. The drawing method is a quick and low-cost way to manufacture inner conductors. The drawing method enables flanges and projections for the inner conductors to be manufactured that are all integrated in the same uniform material piece. Therefore, the inner conductor is mechanically strong.\nSince the inner conductor is deep-drawn, the surface of the inner conductor is extremely smooth, which enables the inner conductor to be readily coated e.g. with silver. Thanks to the smoothness of the surface, the surface area to be coated is smaller than it would be if the surface was uneven. It thus takes less coating material to coat an even surface than an uneven one.\nAn inner conductor manufactured by utilizing the deep drawing method has a small surface resistance, so the electric loss of the resonator remains small and the Q factor of the resonator can be retained good.\nA further advantage of the deep drawing method is that the inner conductor can be manufactured e.g. from a copper blank, in which case the resulting inner conductor does not necessarily have to be coated. The inner conductor manufactured of copper is attached by a specific sleeve, which means that the inner conductor made of copper does not have to be mechanized for the screw threads in a fixing screw.\nSince it is possible to attach the inner conductor by a sleeve, the thickness of the walls of the inner conductor can be retained small, which gives a lightweight inner conductor. The advantage provided by the inner conductor being light is that it is highly tolerant e.g. of vibration. Consequently, external vibration does not easily cause the inner conductor to move or become detached. The structure and attachment of the inner conductor thus enable intermodulation noise to be reduced."}
{"text": "The general aviation market has recently been introduced to the “very small” turbofan engine (i.e. 2000 pounds thrust and less). Simply scaling down larger conventional turbofan engines, however, presents difficulties due mainly to the disproportionate scaling of certain factors, such as strength-to-weight and tolerances.\nThe engine case, such as that depicted in FIGS. 1 and 2, is subjected to asymmetric loading relative to the engine mounts, caused by loads exerted through the bearings, such as engine thrust, foreign object impacts and blade-off events, and caused by inertia loads caused by the engine weight which of course must be supported. These asymmetric loads result in bending moments and shears which must be transmitted through the engine case to the engine mounts. The prior art generally relies on thick walled structures, such as cast engine case components (such as 202 204, 206, 208, 210, 211 in FIG. 1), to react these bending moments in plate bending. Plate bending, however, requires thicker-walled casings to resist and carry bending forces without failure. In very small engines, however, thick casing become a significant component in overall engine weight.\nAn alternate approach is shown in U.S. Pat. No. 4,132,069, which provides an integrally-webbed structure for transferring loads, and in particular bending, through an engine and nacelle structure, so that nacelle loads can be passed to the engine. The scheme, however, adds components to the engine, which reduces reliability, and increases weight and cost. Improvement in engine case technology is therefore desired."}
{"text": "Field of the Invention\nThis invention relates to methods and apparatus for modeling, evaluating and simulating hydrocarbon bearing subterranean formations (which are commonly referred to as reservoirs).\nState of the Art\nPetroleum consists of a complex mixture of hydrocarbons of various molecular weights, plus other organic compounds. The exact molecular composition of petroleum varies widely from formation to formation. The proportion of hydrocarbons in the mixture is highly variable and ranges from as much as 97% by weight in the lighter oils to as little as 50% in the heavier oils and bitumens. The hydrocarbons in petroleum are mostly alkanes (linear or branched), cycloalkanes, aromatic hydrocarbons, or more complicated chemicals like asphaltenes. The other organic compounds in petroleum typically contain carbon dioxide (CO2), nitrogen, oxygen and sulfur, and trace amounts of metals such as iron, nickel, copper and vanadium.\nPetroleum is usually characterized by SARA fractionation where asphaltenes are removed by precipitation with a paraffinic solvent and the deasphalted oil separated into saturates, aromatics and resins by chromatographic separation.\nThe saturates include alkanes and cycloalkanes. The alkanes, also known as paraffins, are saturated hydrocarbons with straight or branched chains which contain only carbon and hydrogen and have the general formula CnH2n+2. They generally have from 5 to 40 carbon atoms per molecule, although trace amounts of shorter or longer molecules may be present in the mixture. The alkanes include methane (CH4), ethane (C2H6), propane (C3H8), i-butane (iC4H10), n-butane (nC4H10), i-pentane (iC5H12), n-pentane (nC5H12), hexane (C6H14), heptane (C7H16), octane (C8H18), nonane (C9H20), decane (C10H22), hendecane (C11H24)—also referred to as endecane or undecane, dodecane (C12H26), tridecane (C13H28), tetradecane (C14H30), pentadecane (C15H32) and hexadecane (C16H34). The cycloalkanes, also known as napthenes, are saturated hydrocarbons which have one or more carbon rings to which hydrogen atoms are attached according to the formula Cycloalkanes have similar properties to alkanes but have higher boiling points. The cycloalkanes include cyclopropane (C3H6), cyclobutane (C4H8), cyclopentane (C5H10), cyclohexane (C6H12), cycloheptane (C7H14), etc.\nThe aromatic hydrocarbons are unsaturated hydrocarbons which have one or more planar six-carbon rings called benzene rings, to which hydrogen atoms are attached with the formula CnHn. They tend to burn with a sooty flame, and many have a sweet aroma. Some are carcinogenic. The aromatic hydrocarbons include benzene (C6H6) and derivatives of benzene as well as polyaromatic hydrocarbons.\nResins are the most polar and aromatic species present in the deasphalted oil and, it has been suggested, contribute to the enhanced solubility of asphaltenes in crude oil by solvating the polar and aromatic portions of the asphaltenic molecules and aggregates.\nAsphaltenes are insoluble in n-alkanes (such as n-pentane or n-heptane) and soluble in toluene. The C:H ratio is approximately 1:1.2, depending on the asphaltene source. Unlike most hydrocarbon constituents, asphaltenes typically contain a few percent of other atoms (called heteroatoms), such as sulfur, nitrogen, oxygen, vanadium and nickel. Heavy oils and tar sands contain much higher proportions of asphaltenes than do medium-API oils or light oils. Condensates are virtually devoid of asphaltenes. As far as asphaltene structure is concerned, experts agree that some of the carbon and hydrogen atoms are bound in ring-like, aromatic groups, which also contain the heteroatoms. Alkane chains and cyclic alkanes contain the rest of the carbon and hydrogen atoms and are linked to the ring groups. Within this framework, asphaltenes exhibit a range of molecular weight and composition. Asphaltenes have been shown to have a distribution of molecular weight in the range of 300 to 1400 g/mol with an average of about 750 g/mol. This is compatible with a molecule contained seven or eight fused aromatic rings, and the range accommodates molecules with four to tens rings. It is also known that asphaltene molecules aggregate to form nanoaggregates and clusters.\nThe life cycle of a reservoir typically follows certain stages including, but not limited to exploration, assessment, reservoir development, production, decline, and abandonment of the reservoir. Important decisions must be made at each of these stages in order to properly allocate resources and to assure that the reservoir meets its production potential. In the early stages of the life cycle, one begins with almost complete ignorance about the distribution of the internal properties within the reservoir. As development continues, diverse types of reservoir data are collected, such as seismic, well logs, and production data. Such reservoir data are combined to construct an understanding of the reservoir.\nComputer-based software applications are commercially available for generating geological models which predict and describe the rock properties and features of subterranean formation. For example, geological models are built from data acquired during the exploration stage, such as seismic analysis, formation evaluation logs, and pressure measurements. Fluid models are built with the input from lab pressure-volume-temperature (PVT) analyses, geochemistry studies, pressure gradients, and downhole fluid analysis (DFA). Fluid models can be combined with geological models as part of a reservoir simulation grid (also commonly referred to as a reservoir model). The reservoir simulation grid represents the three-dimensional physical space of the formation by an array of discrete cells, delineated by a grid system which may be regular or irregular. Values for rock properties (e.g., porosity, permeability, water saturation) and fluid properties (e.g., compositions of liquid and gaseous phases, pressure, and temperature) are associated with each cell. Equations and associated computations are used to model and simulate the flow of fluids during production. Uncertainty in the values of the rock and fluid properties of the reservoir can be investigated by constructing several different realizations of the sets of property values. The phrase “reservoir characterization” is sometimes used to refer to reservoir modeling activities up to the point where the reservoir simulation grid characterizes the static rock and fluid properties of the reservoir, i.e., before the simulation of the dynamic flow of fluids during production.\nSuch computer-based reservoir modeling applications are used to achieve a better understanding of the reservoir and make critical decisions with respect to reservoir development. However, prior to the reservoir development stage, the uncertainty in these models is relatively high. Consequently, known reservoir modeling applications are not always available with sufficient accuracy to permit efficient reservoir development. This is a problem because relatively greater risk exists in the reservoir development stage in comparison with the exploration and assessments stages. Activity tends to occur at a faster pace in the reservoir development stage. For example, an operator typically decides which zones are to be completed immediately after logging and sampling operations. The zones are selected based on predicted commercial value as indicated by the volume of reserves represented in the existing models. If a mistake is made because of model inaccuracy, a costly workover operation and delayed production may result. The risks are particularly high in the case of offshore development because of higher development and operating costs.\nOne particular impediment to efficient reservoir development is reservoir compartmentalization. Reservoir compartmentalization is the natural occurrence of hydraulically isolated pockets within a reservoir. In order to produce a reservoir in an efficient manner, it is necessary to know the structure of the rock and the level of compartmentalization. A reservoir compartment does not produce unless it is tapped by a well. In order to justify the drilling of a well, the reservoir compartment must be sufficiently large to sustain economic production. Furthermore, in order to achieve efficient recovery, it is generally desirable to know the locations of as many of the reservoir compartments as practical before extensive development has been done.\nThere are three industry standard procedures widely used to understand reservoir compartmentalization. First is the evaluation of petrophysical logs. Petrophysical logs may identify impermeable barriers, and the existence of such barriers can be taken to mean that the reservoir is compartmentalized. Examples include gamma ray and NMR logs, both of which can identify impermeable barriers in favorable situations. Another example is the evaluation of mud filtrate invasion monitored by resistivity logs. However, impermeable barriers may be so thin that they are not observable by these logs, or barriers observed by these logs may not extend away from the wellbore and therefore may not compartmentalize the reservoir. Second is the evaluation of pressure gradients. If two permeable zones are not in pressure communication, they are not in flow communication. However, the presumption that pressure communication implies flow communication has repeatedly been proven to be incorrect. Pressure equilibration requires relatively little fluid flow and can occur more than 5 orders of magnitude faster than fluid compositional equilibration, even in the presence of flow barriers. Continuous pressure gradients are a necessary but insufficient test for reservoir connectivity. Third is the comparison of geochemical fingerprints of fluid samples acquired from different locations in the reservoir. Petroleum is a complex chemical mixture, containing many different chemical compounds; the composition of that petroleum can therefore be treated as a fingerprint. If the composition of petroleum samples from two different places in the reservoir is the same, it is assumed that fluids can flow readily between those two places in the reservoir and hence that the reservoir is connected. However, forces such as biodegradation and water washing can occur to different extents in different parts of the reservoir, causing two places in the reservoir to have different fingerprints even if they are connected. Additionally, petroleum samples generated from the same source rock may have very similar fingerprints even if they come from locations in the reservoir that are presently disconnected.\nAn alternative method to assess connectivity is to evaluate hydrocarbon fluid compositional grading. The chemical composition of petroleum must be different in different parts of a connected reservoir. This change in composition with position (typically depth) in the reservoir is referred to as compositional grading. The magnitude of this compositional grading (i.e., the difference in the composition of two fluids collected from different depths), in connected reservoirs at thermodynamic equilibrium, can be modeled with mathematical equations of state (EOS) and measured with downhole fluid analysis. If the magnitude of compositional grading is measured, and the measurement matches the predictions of the model, then the assumptions of the model are believed to be correct. In this case, the assumptions are that the reservoir is connected and at thermodynamic equilibrium. In the event that the magnitude of the compositional grading does not match the predictions of the EOS model, it can be assumed that there is reservoir compartmentalization or that the reservoir fluids are not in equilibrium. Many different forces can contribute to a lack of thermodynamic equilibrium, such as tar mats, water washing, biodegradation, real-time charging, etc. It can be difficult to determine whether the reservoir is compartmentalized or not in thermodynamic equilibrium, and this determination can be critical to important development decisions. More specifically, the traditional EOS (such as the Peng-Robinson EOS developed in 1976) utilized for compositional grading analysis are derived by adding correction terms to the ideal gas law to address gas-liquid equilibria. Thus, these standard EOS allow for compositional analysis of only gas and liquid phase fractions of the reservoir fluid, and such limited information makes it difficult to determine whether the reservoir is compartmentalized or not in thermodynamic equilibrium.\nThus, there is a clear need for methodologies that provide for an effective understanding of reservoir compartmentalization as early as possible (e.g., before development) in the lifecycle of the reservoir."}
{"text": "A typical multi-color dye donor web that is used in a dye transfer or thermal printer is substantially thin and has a repeating series of three different rectangular-shaped color sections or patches such as a yellow color section, a magenta color section and a cyan color section. In addition, there may be a transparent colorless laminating section immediately after the cyan color section.\nEach color section of the dye donor web consists of a dye transfer area which is used for dye transfer printing and a pair of opposite longitudinal edge areas alongside the dye transfer area which often are not used for printing. The dye transfer area may be about 152 mm wide and the two longitudinal edge areas may each be about 5.5 mm wide, so that the total web width is approximately 163 mm.\nTo make a multi-color image print using a thermal printer, a motorized donor web take-up spool draws a longitudinal portion of the dye donor web off a donor web supply spool in order to successively move an unused single series of yellow, magenta and cyan color sections over a stationary liner array (bead) of selectively heated resistive elements on a thermal print head between the supply and take-up spools. Respective color dyes within the yellow, magenta and cyan color sections are successively heat-transferred line-by-line, via the selectively heated resistive elements, onto a dye receiver medium such as a paper or transparency sheet or roll, to form the color image print. The selectively heated resistive elements often extend across the entire width of a color section, i.e. across the dye transfer area and the two longitudinal edge areas comprising that color section. However, only those resistive elements that contact the dye transfer area are selectively heated. Those resistive elements that contact the two longitudinal edge areas are not heated. Consequently, the dye transfer occurs from the dye transfer area to the dye receiver medium, but not from the two longitudinal edge areas to the dye receiver medium.\nAs each color section is drawn over the selectively heated resistive elements, it is subjected to a longitudinal tension particularly by the forward pulling force of the motorized donor web take-up spool. Since the dye transfer area in the color section is heated by the resistive elements, but the two longitudinal edge areas alongside the dye transfer area are not, the dye transfer area is significantly weakened and therefore vulnerable to stretching as compared to the two longitudinal edge areas. Consequently, the longitudinal tension will stretch the dye transfer area relative to the two longitudinal edge areas. This stretching causes the dye transfer area to become thinner than the non-stretched edge areas, which in turn causes some creases or wrinkles to develop in the dye transfer area, most acutely in those regions of the dye transfer area that are close to the non-stretched longitudinal edge areas. The creases or wrinkles occur most acutely in the regions of the dye transfer area that are close to the non-stretched edge areas because of the sharp, i.e. abrupt, transition between the stretched (thinner) transfer area and the non-stretched (thicker) edge areas.\nAs the dye donor web is pulled by the motorized donor web take-up spool over the selectively heated resistive elements, the creases or wrinkles tend to spread from a trailing (rear) end portion of a used dye transfer area at least to a leading (front) end portion of the next dye transfer area to be used. A known problem that can result is that the creases in the leading (front) end portion of the next dye transfer area to be used will cause undesirable line artifacts to be printed on a leading (front) end portion of the dye receiver medium. The line artifacts printed on the dye receiver medium, although they may be relatively short, are quite visible.\nThe question presented therefore is how to solve the problem of the creases or wrinkles being created in an unused transfer area so that no line artifacts are printed on the dye receiver medium during the dye transfer.\nThe Cross-Referenced Applications and Patent\nThe cross-referenced applications each disclose a thermal printer capable of preventing crease formation in a dye transfer area of a dye donor web that can cause line artifacts to be printed on a dye receiver during the dye transfer from the dye transfer area to the dye receiver. To prevent crease formation, there is provided a crease-preventing platen roller that is movable to hold a dye transfer area and the two longitudinal edge areas alongside the dye transfer area against a print head. The crease-preventing roller has a pair of roller end portions that apply a constant pressure against the two longitudinal edge areas, and a roller main portion between the roller end portions that applies a lesser pressure against the dye transfer area. Since the pressure applied against the two longitudinal edge areas is greater than the pressure applied against the dye transfer area, the two longitudinal edge areas will be stretched the same as the dye transfer area, so that creases will not be formed in the dye transfer area. This is so notwithstanding that the dye transfer area is heated by the print head, but the two longitudinal edge areas are not.\nIn contrast to the cross-referenced applications, the referenced incorporated (prior art) patent discloses a thermal printer that is adapted to optimize print image quality by preventing undesired pressure variations along the line of contact between the dye donor web and the linear array of selectively heated resistive elements in the thermal print head (the patent does not discuss the problem of crease formation). To optimize print image quality, there is provided a plurality of pressure applying rods that bear down on the thermal print head to urge the selected heated resistive elements into pressure contact with the dye donor web. The amount of pressure applied by each rod and the location of each rod along the print head is individually adjusted in response to sensed changes in different operating parameters that negatively affect print image quality, such as print head temperature when printing dark vs. light image portions, and thickness and/or stiffness of the dye donor web."}
{"text": "The field of the invention is submarine pipeline trenching apparatus.\nApparatus is presently known for digging a trench below a pipeline, examples of which are U.S. Pat. Nos. 3,429,131 and 3,429,132 granted to Charles F. Martin. With such prior patents it was necesary to pre-form or pre-dig a hole below the pipeline in the ocean bottom for receiving the cutter assembly or assemblies so that the apparatus could be mounted on the pipeline in its operable digging position. Additionally, difficulty was sometimes experienced in use due to shifting and other instability of the apparatus as it moved along the pipeline because of lateral bends and vertical undulations. Dragging of the lower ends of the cutters along the bottom of the trench often interfered with the cutting action and rapid longitudinal travel of the cutters relative to the pipeline."}
{"text": "1. Field of the Invention\nThe present invention relates to a process and apparatus plant for treating food products with ozone, aiming especially to bleach or decolorize, sanitize and deodorize the products thus treated.\n2. Description of the Related Art\nThe literature relating to the ozone treatment of food products, in particular in the field of seafood (fish, crustacea, etc.) is known to be very extensive, and the reader may be referred for instance to documents FR-385,815, EP-294,502, FR-797,928 or else U.S. Pat. No. 4,559,902.\nThe use of ozone has thus been most particularly described in the case of the treatment of fish meat with the objective of sterilizing and deodorizing the meat, the treated fish meat coming from the recovery of the residues remaining on the bones and the head after filleting the fish, from fillet cutting scraps or even from actual fillets, the meat in question being eventually used as raw material for the manufacture of various products, such as surimi seafood products or other fish pastes, terrines or steaks.\nA typical apparatus for the recovery and treatment of such fish meat for producing what those skilled in the art call \"base surimi\" seafood therefore comprises the following steps:\na step of separating the meat from the bones and heads (resulting in a first crude pulp, which is ready to use if necessary); PA1 one or more operations of washing in water which optionally is slightly acidified, each washing operation being followed by a draining step (resulting in a washed pulp ready to use for some applications); PA1 a refining step making it possible, by passing the material over a screen, to separate the proteins from the impurities (skin residues, etc.); PA1 a final operation of mechanical water/meat separation carried out by centrifugal sedimentation or screw pressing, resulting in washed and refined fish meat ready to use (base surimi seafood). PA1 to improve the productivity of ozone-treated food products; PA1 to improve the quality of the transfer of ozone to the food product to be treated (it is known in practice that at the present time only 15 to 60%--depending on the plants available--of the ozone injected is actually transferred to the food product in question); PA1 to avoid, nevertheless, the risk of impairing the product (by way of illustration, mention may be made here of the risk of impairing fish meat by turning it brown because of local overdosing); PA1 to reduce overall the number of washing steps carried out in the user treatment chain; PA1 to achieve higher quality in the product. PA1 a) a supply of the initial solution containing the product is used; PA1 b) a pumping device is used which allows the initial solution to be taken under pressure to a contactor; PA1 c) an ozone-containing treatment gas mixture is injected into the initial solution, the injection taking place at one or simultaneously at several of the following locations: PA1 the ozone dose used for the treatment, expressed in grams of ozone per kilo of treated product, lies within the range going from 0.2 to 2 g/kg (taking into account not only the specificity of the treated product and of the desired treatment, but also, for example, the national legislation governing each user site). PA1 The ozone content in the treatment gas mixture lies within the range going from 10 to 200 g/m.sup.3, preferably lies within the range going from 20 to 120 g/m.sup.3 and more preferably within the range going from 40 to 100 g/m.sup.3 of mixture. PA1 The composition of the initial solution treated satisfies a degree of dilution corresponding to 1 volume of product per 0.5 to 10 volumes of water, but preferably to a degree of dilution corresponding to 1 volume of product per 1 to 5 volumes of water and even more preferably to a degree of dilution corresponding to 1 volume of product per 3 to 5 volumes of water. PA1 a) a supply of the initial solution to be treated; PA1 b) a pumping device suitable for the initial solution to be taken under pressure to a contactor; PA1 c) a supply of a treatment gas mixture which contains ozone; PA1 d) means for injecting the treatment gas mixture into the solution, allowing the mixture to be injected at one or simultaneously at several of the following locations:\nThe results of animal meat bleaching are commonly monitored by industrial sites by color measurements, called colorimetry, with readings of the brightness or whiteness, of the red/blue index and the yellow/green index, for example using the conventional L*/a*/b* system (CIE 1976 reference), the factor L being expressed in %.\nFor reasons of simplification, throughout the following text reference will be made, indiscriminately, to the L/a/b or L,a,b system or to their results, clearly keeping in mind that such references pertain to the abovementioned evaluation system.\nThe desired objectives, in terms of colorimetry, throughout the various washing operations carried out, are generally (depending on the product in question) an increase in the whiteness, a decrease in the redness and the absence of any change in yellows/browns or even a reduction in the latter.\nThe aforementioned document EP-A-284,502 describes in particular an apparatus for treating animal meat, consisting of a shell in which a hollow tube provided with a helical partition rotates, the whole assembly constituting a transverse screw, the free volume existing between the internal walls of the outer shell and the external helical partition of the hollow tube constituting a splashing chamber through which the meat/water mixture to be treated travels, the gas mixture containing ozone being injected into the hollow tube and diffusing outward into the splashing chamber through a plurality of diffusers located along the transfer screw.\nThe example provided by the document indicates a meat transfer rate inside the system of about 1 cm/s, which gives, taking into account the geometrical characteristics of the apparatus, a meat output of about 1 liter/minute.\nApart from the complexity of the system described in that document, work successfully carried out by the Applicants has allowed them to show that, because of the dynamics created in such an apparatus, the performance in terms of contact between the fish meat and the dissolved ozone is insufficient (no intimate mixing between the meat and the gas), leading to solid/gas demixing, but also, as a consequence, a lack of effectiveness over the entire diameter of the screw.\nIt may also be added that the use of gas injection pores for injecting the gas seems not to be very compatible with the cleanability requirements commonly practiced in the food industry (zones of the screw that are difficult to reach by a cleaning agent and potential blocking of the porous injectors with food material)."}
{"text": "This invention is concerned with production of a simple, economical carton for bottles, particularly designed for capped bottles having sloping necks leading upwardly from shoulder at the top of a cylindrical lower bottle portion. Such bottles are used in the brewing and soft drink beverage industries. The carton design results in considerable material savings in comparison to a normal rectangular carton, and assures tight packing of the bottles against one another and against the carton walls. This eliminates the necessity of an interior partition, which can be used optionally when desired. The carton also provides multi-layer cushioning about the bottles for external protection.\nThe usual bottle carton has a rectangular configuration and is made from paperboard or corrugated paperboard. In order to provide clearance for normal top-loading of the carton, it is necessary to provide partitions or bottle carriers with dividers between the adjacent bottles. The interior area of such a carton contains considerable unused space, particularly along the upper portions of the bottles where the rectangular corners are spaced at a distance outward from the sloping bottle necks.\nPrior bottle carriers and packages have tackled the problem of this unused space, but this has been done primarily with \"wrap-around\" carriers which are fitted about the bottles in a single operation and which provide only limited bottle protection at the ends of the carrier. Patents showing this general feature are too numerous to fully list, but examples are shown in U.S. Pat. No. 3,186,545 to Conrades; 3,167,214 and 3,152,688 to Mahon; and 3,306,519 to Wood. These patents include cutouts for the caps of the bottles and tabs or dividers to separate individual bottles from one another.\nU.S. Pat. No. 3,578,238 shows a rather tight wrap-around carton without cutouts, again formed in a single wrapping operation. This carton is tapered upwardly on all four sides to firmly grasp the bottles. The container shown is a lightweight carrier and not designed for use as a shipping carton. A related carrier construction is shown in U.S. Pat. No. 3,688,972 to Mahon.\nU.S. Pat. No. 3,519,127 to Wood shows a carrier open at both ends, and including yieldable apertures in a top panel. The carrier is fitted over the bottles to draw their upper ends together due to the fact that the apertures are spaced apart a distance less than the bottle centers. As the carrier is completed at the bottom of the bottles, the bottles are pulled into a parallel relationship. The apertures are formed with slits along the sides of the carrier to yield as the necks of the bottles move apart.\nThe carton described herein is designed for shipping of bottles, but is also capable of use as a small carrier package for six or twelve bottles. It is designed to package two or more rows of upright capped bottles. It can be manufactured from relatively stiff corrugated board products. The described method of assembly of the package results in a very tight package with sufficient yieldability and cushioning of the bottles to assure safe shipment without necessitating the use of interior partitions or other dividers between adjacent bottles."}
{"text": "1. Field of the Invention\nThe present invention relates to a document management apparatus and a document management method and a computer-readable storage medium storing a computer program. More particularly, the present invention relates to a document management apparatus having a construction capable of proving objective validity of document data, managed in a database, microfilm or the like, i.e., the date and time, nonoccurrence of tampering of the document data and the like, and a document management method and a computer program.\n2. Description of the Related Art\nIn a construction to store and manage various document data including text data, image data and program data, when document data is stored in a database connected to e.g. a network, there is a probability of data tampering or the like by external unauthorized access. Recent document management systems employ various arrangements to prove the date and time of generation, nonoccurrence of tampering and the like of management subject document data.\nIt is known to provide a technique for proving the date and time of generation, nonoccurrence of tampering and the like of management subject document data, an arrangement utilizing time stamping. Time stamp is proof data issued from a time stamping authority (TSA) as a trusted third party. A digital signature as encryption data based on a private key from the time stamping authority is generated for data obtained by combining a hash value based on document data provided by a document management system with reliable time information obtained by the time stamping authority (TSA), and issued to the document management system. The protocol utilizing a time stamp issued from the time stamping authority (TSA) is referred to as a “simple protocol”.\nThe simple protocol will be described with reference to FIG. 1. As shown in FIG. 1, a document management apparatus 11 which manages various documents transmits a hash value (H) based on management subject document data to a time stamping authority (TSA) 12. The time stamping authority (TSA) 12 generates a digital signature by using a private key from the time stamping authority (TSA) 12 compliant with Public Key Infrastructure (PKI) for data obtained by combining the received hash value (H) with time information (Time) received from a reliable time information supply source 13. That is, a digital signature (Sig.) is applied to [hash value (H)+time information (Time)], thereby proof data 21 as encrypted data is generated, and transmitted to the document management apparatus 11.\nThe document management apparatus 11 links the proof data 21, received from the time stamping authority (TSA) 12, with the document data to be stored in the database, i.e., document data as the source of the hash value (H), and stores the proof data 21.\nVerification processing based on the proof data 21 is verification processing of the digital signature of the proof data 21, i.e., decrypting of the digital signature based on a public key from the time stamping authority (TSA) 12. The data as the combination of the hash value (H) and the time information (Time) is obtained by decrypting the digital signature by using the public key from the time stamping authority (TSA) 12. Further, the date and time of generation, nonoccurrence of tampering and the like of the document are verified by checking the correspondence between the hash value (H) obtained by decrypting the proof data 21 and a hash value (H′) newly generated based on the document data managed in the document management apparatus 11. In this simple protocol, a public key and a private key compliant with the Public Key Infrastructure (PKI) are employed, and the public key applied upon verification processing based on the proof data 21, i.e., the public key from the time stamping authority (TSA) 12 is obtained from a public key certificate issued from a Certificate Authority (CA), thereby creditability of public key is guaranteed. However, generally, a validity period is set in a public key certificate since it can be considered that permanent use of the same key is prevented and the safety can be increased by updating a pair of public and private keys periodically.\nHowever, when such key update is performed, update of the digital signature in the proof data 21 is required. That is, to generate a digital signature and perform verification using an updated valid pair of public and private keys, the proof data is updated upon each key update processing. As described above, the simple protocol has a problem that update processing is required due to the use of an encryption key based on the Public Key Infrastructure (PKI). Next, a linking protocol as a data proof protocol without Public Key Infrastructure (PKI) will be described with reference to FIG. 2.\nAs shown in FIG. 2, in the linking protocol, a hash value transmitted from a document management apparatus 31 has a sequence number n (1, 2, 3 . . . ). A time stamping authority (TSA) 32 generates document authentication linking data [Ln] as a new hash value based on an n-th hash value [Hn] received from the document management apparatus 31, based on document authentication linking data [Ln−1] as a hash value generated based on a hash value previously received from the document management apparatus 31, the sequence number [n] and the newly received hash value [Hn]. That is,Ln=Hash(Ln−1, n, Hn)The new document authentication linking data [Ln] is generated based on the above expression.\nThe time stamping authority (TSA) 32 generates the next document authentication linking data [Ln] based on immediately-previously generated document authentication linking data [Ln−1], and publicizes the linking data at a predetermined timing in a publication indicating a certifiable date of issuance such as a newspaper. In FIG. 2, linking data [LM] 41 is publicized in a newspaper 51 issued on Oct. 20, 2004, and linking data [LN] 42 is publicized in a newspaper 52 issued on Apr. 12, 2005.\nThe time stamping authority (TSA) 32 stores all the hash values [Hn] sequentially received from the document management apparatus 31, all the linking data [Ln] and sequence data [n].\nIn this manner, as the time stamping authority (TSA) 32 performs hash computation using immediately previous linking data and newly-received hash value upon linking data [Ln] generation processing, the context of all the linking data can be specified. Verification processing is performed as follows. When data verification is performed upon publicizing of linking data (Ln), the document management apparatus generates a hash value (HN′) based on document data, and calculates a hash value [LN′] based on the generated hash value (HN′), linking data (LN−1) held in the time stamping authority (TSA) 32 and a sequence number [N]. That is, the hash value [LN′] is calculated asLN′=Hash(LN−1, N, HN′).If the calculated hash value [LN′] corresponds with the publicized linking data (LN) in the newspaper, nonoccurrence of tampering is proved.\nIn the linking protocol, the proof data update accompanying the key data update compliant with the Public Key Infrastructure (PKI) as in the above-described simple protocol is unnecessary. However, as the time stamping authority (TSA) 32 needs to hold all the hash values [Hn] sequentially received from the document management apparatus 31, all the linking data [Ln] and sequence data [n] , the load on data management is heavy. Further, as the linking data is periodically publicized in publication such as a newspaper, the cost for publicizing is increased."}
{"text": "1. Field of the Invention\nThis invention relates to computer-based communication, and more particularly to a method, system and program for switching communication techniques during the course of a communications session.\n2. Description of the Related Art\nThe following descriptions and examples are not admitted to be prior art by virtue of their inclusion within this section.\nThe continuing proliferation of powerful, convenient computational devices has been accompanied by an increase in the use of networks connecting these devices. Computational devices include computers and other, often portable, devices such as wireless telephones, personal digital assistants (PDA's), automobile-based computers and appliance-based computers. Such computational devices are also sometimes termed “pervasive devices.” “Computer,” as used herein, may refer to any of such computational devices. The networks connecting computational devices may be “wired” networks, formed using “land lines” such as copper wire or fiber optic cable, wireless networks employing earth and/or satellite-based wireless transmission links, or combinations of wired and wireless network portions. Networked computational devices are increasingly configured with the ability to interact and communicate with one another. For example, electronic mail may currently be sent to some wireless telephones and pagers, and PDA's may be used to access some Internet and World Wide Web (or “web”) sites.\nThe many available techniques for computer-based communication exhibit a range of values for criteria such as immediacy of communication, security, and bandwidth required. Email, for example, is analogous to a letter or memo, with a low degree of immediacy, relatively low bandwidth requirement, and potentially high security. AnIP (Internet Protocol) conference, or meeting held over the Internet using cameras, by contrast, has high immediacy, high bandwidth requirements, and potentially lower security (since encryption measures may detract from the ability to provide real-time communication). Other computer-based communication techniques (also referred to herein as “communications applications”), such as instant messaging, message boards, internet chat, and IP telephone, can be described using similar criteria.\nThis wide array of communication options can cause inconvenience and difficulty for user's in some cases. For example, it may be difficult to determine the best communications application to use for initiating communications with a particular person. This may occur when the availability of that person using various communications applications is not known. “Availability” as used herein may refer to the availability of a person (e.g., is the person logged onto his computer), the availability of equipment (such as a web camera), and/or the availability of needed software. In other words, availability may encompass numerous circumstances that might limit the ability to contact a person by a particular communications technique. Furthermore, the identifier of the person within a desired communications application, commonly called a “user id”, may not be known by the person wishing to initiate communication. Before a communication session can be initiated, therefore, a “pre-conversation” by telephone or some other method may be needed in order to ascertain availability and/or obtain the person's identifier. Such additional conversations can greatly increase the time and effort expended in carrying out computer-based communication, particularly when more that two people are involved in the communication.\nEven when a communications session is initiated using a desired communications application, it may be realized during the course of the conversation that another form of communication would be more convenient. This could happen if the nature of the conversation changes, or if it simply becomes clear that the best communications application wasn't used from the beginning. Changing to another communications application can be rather inconvenient, however. Each participant may have to open up another application program, and efforts as described above to verify availability and obtain participant identifiers may need to be carried out. These steps necessitate enough disruption to a flow of conversation that typically changing to another communications application is avoided. Instead, participants may continue to “stretch” the capabilities of a given application. For example, participants may continue to send repeated short email messages back and forth, even though many “mouse clicks” could be saved if a messaging or chat application were used instead.\nIt would therefore be desirable to develop a way to make it convenient for participants in computer-based communications to select the most appropriate communications application."}
{"text": "Applanation tonometers.\nThis invention improves on the tonometer of U.S. Pat. No. 5,070,875, the disclosure of which is hereby incorporated by reference. The \"\"875 patent discloses an applanation tonometer having a replaceable prism; and this invention improves on the prism and its replacement, while otherwise relying on-tonometer components such as described in the \"\"875 patent.\nPrisms for applanation tonometers have proved to be problematic. The art has suggested several prism variations for applanation tonometers, including U.S. Pat. Nos. 5,190,042 and 5,203,331. None of these have yet demonstrated operability sufficient to be successful in the marketplace. The problems involve reliable production of an accurate applanation signal and the difficulty of combining this with a tonometer having a conveniently replaceable prism producing reliable operation at a low cost.\nThe improved prism for our tonometer solves these problems partly by arranging emitter and detector ports opposite each other on opposite sides of a longitudinal or viewing axis so that axes of the emitter and detector ports are angled about 45xc2x0 to the longitudinal axis. This makes the light-transmitting surfaces of the emitter and detector ports approximately perpendicular to each other. Internal-reflecting prism surfaces for the emitter and detector are also arranged opposite each other on opposite sides of the longitudinal axis and proximate to the applanation surface. This makes the emitter- and detector-reflecting surfaces approximately parallel to each other. A light path proceeding from a source through the emitter port is angled at about 45xc2x0 to the longitudinal axis and directed across the longitudinal axis for incidence on the emitter-reflecting surface. From there, the light is incident on an applanation surface, and a portion of the light reflected from the applanation surface is incident on the detector-reflecting surface. A light path to the detector port for the light reflected from the applanation surface and the detector-reflecting surface proceeds at an angle of about 45xc2x0 to the longitudinal prism axis across the longitudinal axis and is output at the detector port.\nA simplification of such a prism eliminates the internally reflecting surfaces for the emitter and detector so that light passes directly from the emitter through the emitter port to be incident on the applanation surface from which the light is reflected through the detector port to the detector. Light entering one of the ports normal to the port interface surface is incident on the applanation surface at approximately 45xc2x0 and is reflected to be approximately normal to the opposite port interface surface. Such a prism can have either an elongated or shortened viewing axis, depending partly on the desired configuration of a holder for the prism.\nWith a prism of such a configuration receiving light from a light-emitting diode, which is presently preferred as a source of prism illumination, the light incident on the applanation surface can be made to vary in intensity radially from a center of the applanation surface to a perimeter of the applanation surface. A consequence of such a varying light distribution at the applanation surface is that the amount of reflected light reaching the detector upon concentric contact of the applanation surface with a cornea differs distinguishably from an amount of reflected light reaching the detector upon eccentric contact of the applanation surface with a cornea. The readily distinguishable difference in detected light from concentric and eccentric contact with a cornea is preferably used to determine that the prism is properly positioned to proceed with an intraocular pressure measurement sequence.\nA prism having desirable light-transmitting characteristics is also preferably replaceable and disposable. This is accomplished by provisions for easy insertion and withdrawal of the prism, making the prism inexpensive, providing a system for requiring and detecting replacement of the prism, and making the tonometer quickly and reliably calibrated to each new prism. We also prefer making the tonometer responsive to input for adjusting an intraocular pressure measurement to accommodate variations in corneal curvature and thickness of a particular eye being examined."}
{"text": "The statements in this section merely provide background information related to the present disclosure and may or may not constitute prior art.\nA continuously variable transmission (“CVT”) typically includes gearing that operatively couples a variator between a rotary power source, such as an engine or electric motor, and a final drive unit. The variator includes a rotary input disk and a rotary output disk which are able to steplessly or continuously vary the ratio of an input speed to an output speed (the “variator ratio”). The overall speed ratio provided by the CVT is a function of the variator ratio and the associated gearing. The output disc includes external gear teeth that are in mesh with and drive a corresponding gear. The gear in turn is functionally coupled to an output shaft or layshaft that is functionally coupled to the final drive unit.\nTransmission axial length and mass significantly impact its power density and efficiency. Accordingly, there is a constant need for improved CVT designs that minimize axial length and mass while providing sufficient performance characteristics."}
{"text": "Position recognition systems are used to determine the position of and track a particular object in 3-dimensions. In robot assisted surgeries, for example, certain objects, such as a surgical instrument, needs to be tracked with a high degree of precision as the instrument is being positioned and moved by a physician.\nAn infrared signal based position recognition systems use either passive sensors or active sensors. In passive sensors, objects to be tracked such as spherical balls are positioned at strategic locations of the object to be tracked. Infrared transmitters transmit a signal having a predetermined waveform pattern. The spherical balls reflect the signal and the reflected signals are received by infrared signal detectors that are generally located near the transmitters.\nIn active sensors, the objects to be tracked have their own infrared transmitters and thus generate their own infrared signals in response to the signals transmitted from an infrared transmitter of the position recognition system.\nIn either active or passive sensors, the responsive infrared signals from the objects to be tracked are received and amplified in a form which is as close to the original infrared signals as possible. The amplified signal is then converted into a digital signal which is usable by a signal processor. The system then geometrically resolves the 3-dimensional position of the spherical balls based on one or more of the converted digital signals, cameras, digital waveform shape of the amplified signals, known locations of the spherical balls, distance, the time it took to receive the responsive signals, or a combination thereof.\nOne problem is that due to attenuation of infrared signals, the conventional position recognition systems have a relatively limited working range of distance between the system and objects to be tracked. Typically, the working range is limited to about 3-5 meters, at best, before the amplified signal loses sufficient original waveform information to be useful.\nMoreover, conventional systems require relatively high bandwidth, high slew rate amplifiers which involve complex and expensive circuits. Therefore there is a need to provide an improved system and method for recognizing the 3-dimensional position of an object, which has a longer working range, is accurate, and is less expensive."}
{"text": "Carboxamidine intermediates can be used to prepare oxadiazole compounds which are inhibitors of five lipoxygenase activating protein (FLAP) and are thus useful for treating a variety of diseases and disorders that are mediated or sustained through the activity of leukotrienes including asthma, allergy, rheumatoid arthritis, multiple sclerosis, inflammatory pain, acute chest syndrome and cardiovascular diseases including atherosclerosis, myocardial infarction and stroke. The preparation of oxadiazole compounds via a carboxamidine intermediate is disclosed in U.S. application Ser. No. 13/208,582 “Oxadiazole Inhibitors of Leukotriene Production”, filed Aug. 12, 2011."}
{"text": "This invention relates to a high-speed, high-efficiency image recording apparatus and system, for recording images o photo-sensitive material, and to an image recording roll magazine loaded in an image recording apparatus for recording a color image on photo-sensitive material.\nWhen color images have been recorded by optically applying electrical image signals to a photo-sensitive material, a so-called \"drum type image recording system\" has been employed. In this image recording method, a drum on which a photo-sensitive material of a predetermined size has been wound is rotated, and the photo-sensitive material is scanned by a modulated point light source in a direction parallel to the drum.\nHowever, since the speed of rotation of the drum has an upper limit, it is rather difficult to improve the image recording speed. Furthermore, it takes a relatively long time to attach the photo-sensitive material to the drum or to detach it from the drum, and the mechanism for doing so is relatively intricate. Thus, the image recording apparatus is necessarily bulky and complicated. In addition, the photo-sensitive materials are of a predetermined size and are processed batchwise one at a time, so that the number of photo-sensitive materials which could be processed per unit time is relatively small.\nOne example of the above-described image recording apparatus (Japanese Patent Application Laid-Open No. 279,973/87) now will be described with reference to FIGS. 4 and 5. As shown in FIG. 4, a thermal developing photo-sensitive material 16 is accommodated in a magazine 14 which is loaded in a housing 12. The photo-sensitive material 16 is pulled out of the magazine 14 and then is cut to a predetermined length by a cutter 18. The cut photo-sensitive material 16 is wound around an exposure drum 20 (in the direction of the arrow A). An exposure head 22 having three color point light sources made up of light emitting elements faces the exposure drum with the material 16 disposed therebetween. After one optical image exposure, the exposure drum 20 is turned in the opposite direction, so that the photo-sensitive material 16 comes off the drum by a scraper 24. The material then is conveyed through a water applying unit 26 to a developing and transferring unit 28.\nImage receiving sheets 32 stacked on a tray 30 are delivered to the developing and transferring unit 28 one at a time. In the unit 28, each image receiving sheet 32 is combined with a corresponding piece of photo-sensitive material 16, and then the image receiving sheet 32 and the photo-sensitive material 16 are heated by a heater in the unit 28. Thus, the photo-sensitive material is developed, and the developed image is transferred onto the image receiving sheet 32. Thereafter the photo-sensitive material 16 and the image receiving sheet 32 are separated from each other by a scraper 36, and the photo-sensitive material 16 is delivered into a disposal tray 38, while the image receiving sheet 32 is conveyed through a drying unit 40 to a take-out tray 42.\nAs shown in FIG. 5, the photo-sensitive material 16 wound on the exposure drum 20 is rotated around the axis (in the direction of the arrow A), thus permitting the main scanning operation of the exposure head 22. The exposure head 22 has three primary color light-emitting elements 44, and is moved along a guide 46 arranged parallel to the axis of the exposure drum 20 to perform an auxiliary scanning operation in the direction of the arrow B. In order to move the exposure head 22 in the direction of the arrow B, a pair of pulleys 48 and 50 are arranged in a line which also is parallel to the axis of the exposure drum and a wire 52 is laid over the pulleys 48 and 50, with one end secured to one side of the exposure head 22 and with the other end secured through a tension means 54 to the other side of the exposure head 22.\nThe photo-sensitive material 16 cut to the predetermined length by the cutter 18 is wound on the exposure drum 20, and the drum 20 is rotated in the direction of the arrow A by an electric motor (not shown). As the pulley 48 is turned, the exposure head 22 is moved in the direction of the arrow B to carry out the auxiliary scanning operation. In the auxiliary scanning operation, the light emitting elements 44 perform exposure only when the pulley 48 is driven clockwise (in the direction of the arrow C).\nAfter exposure, referring again to FIG. 4, the exposure drum 20 is turned in the opposite direction, and the photo-sensitive material 16 is removed from the exposure drum 20 by the scraper 24 and then is delivered through the water applying unit 26 to the developing and transferring unit 28. The image receiving sheet 32 taken out of the tray 30 is delivered to the developing and transferring unit 28, where it is combined with the photosensitive material 16 such that the emulsion surfaces of the image receiving sheet 32 and the photo-sensitive material 16 are kept in close contact with each other. Under this condition, the heater in the developing and transferring unit 28 is energized, as a result of which the latent image on the photo-sensitive material 16 is thermally developed, and the image thus developed is transferred onto the image receiving sheet 32. In this case, the image is positively transferred onto the image receiving sheet 32 because water has been applied to the emulsion surface of the thermal developing photo-sensitive material 16 by the water applying unit 26. After the image transferring operation, the photo-sensitive material 16 is removed from the image receiving sheet 32 by the scraper 36 and is then delivered into the disposal tray 38, while the image receiving sheet 32 is delivered into the take-out tray 42.\nAs described above, in an image recording apparatus using thermal developing photo-sensitive material, the photo-sensitive material 16 cut to a certain length is subjected to optical exposure and other treatments after being wound on the exposure drum 20. The operations of cutting the photo-sensitive material, of attaching the photo-sensitive material to the exposure drum, and of removing it from the exposure drum take a relatively long period of time, and require an intricate mechanism. Furthermore, in the image recording apparatus, an image is recorded on the photo-sensitive material 16 by the main and auxiliary scanning operations which are effected with the three light emitting elements 44 provided on the exposure head 22 to give exposure in three colors according to the photo-sensitivity characteristic of the photo-sensitive material 16, and therefore the image recording time is considerably long. In addition the image recording size is limited by the diameter of the exposure drum 20. Furthermore, since it is necessary to provide the scraper 36 for removing the image receiving sheet 32 from the photo-sensitive material 16, and the disposal tray 38 for receiving the photo-sensitive material 16 used, the apparatus is unavoidably bulky."}
{"text": "Embodiments of the present specification generally relate to a power generation system having an energy storage system, and more particularly to an energy storage system having an auxiliary direct current (DC) bus to provide power to auxiliary loads using multiple energy sources.\nAs will be appreciated, an energy storage system is an electrical power storage and distribution system. Typically, the energy storage system includes power buses, energy storage devices, power electronic circuits, and auxiliary loads. The power electronic circuits include inverters and converters.\nTypically, the power generation system and in particular, the energy storage system includes a main bus. The auxiliary loads are coupled to the main bus. Traditionally, a single energy source is used to supply power to the auxiliary loads via the main bus. Any fault in the energy source results in a disruption in supply of power to the auxiliary loads, thereby compromising the reliability of the power generation system. Also, use of a single DC bus to provide power to the auxiliary loads may result in a complex architecture of the power generation system.\nIn recent times, in addition to the main bus, the power generation system includes additional buses to supply power to the auxiliary loads. Even though the additional buses are employed, faults in the energy source that supplies power to the auxiliary loads tend to lower the reliability of the associated power generation system."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of panels for forming walls of buildings, and more specifically to a prefabricated panel having a concrete planar portion with at least one concrete rib extending from a face thereof and at least one block of polystyrene fitted adjacent thereto, and a metal stud channel containing a polystyrene strip secured over each at least one rib for receiving dry wall screws. The present invention also relates to a structure formed of such panels.\n2. Background Information\nConcrete block and other prior art wall designs formed upright and incrementally generally have irregular faces and weak cohesion. The placing of concrete blocks is only as accurate as the skill of the blocklayer permits. The resulting exposed, exterior face of the wall is always uneven to some extent, requiring a substantial stucco covering. More important is that the mortar used to join the concrete blocks is below the strength of properly hardened concrete and forms a relatively weak bond with the blocks. Hurricanes and other natural disasters can sometimes shatter such walls. Finally, the insulating properties of such walls are at most adequate.\nAttempts have been made to solve these problems with a preformed wall panel. The panel of Walston, U.S. Pat. No. 4,512,126, issued on Apr. 23, 1985, for example, is formed of two layers of polystyrene fitted into a ribbed concrete shell in a multiple step, multiple pour process. Anchor hooks must be embedded in the floor slab when it is poured. Problems with Walston include time-consuming and expensive forming methods and awkwardness and complexity in anchoring procedures.\n3. Description of the Prior Art\nThere have long been prefabricated wall panels for rapid construction of buildings. None teach panels combining high strength, maximum insulating properties, single pour formation, and solid, easy to use anchoring means.\nAdditional examples of prior panels include Lewis, U.S. Pat. No. 4,494,353, issued on Jan. 22, 1985. Lewis teaches a rigid, rectangular insulating member which fits within the U of opposing vertical hat channels and a bottom channel. The bottom channel is welded to a bottom plate which is secured by anchor bolts to a concrete foundation. Panel connector strips are welded to the hat channels and reinforcing rods extend through holes in the connector strips. Gunite is then sprayed over this structure to form the finished panel. Numerous and complicated metal parts make Lewis expensive. Also, two separate layers of concrete are required, adding pouring and curing time and expense.\nKinard, U.S. Pat. No. 4,532,745, issued on Aug. 6, 1985, discloses a wall formed of foam blocks and channels. The blocks, which have periodic vertical bores, are placed edge to edge. Then a channel with periodic holes is placed along the common top edge of the blocks so that the holes align with the vertical bores. Concrete is poured through the holes into the vertical bores, filling the bores and the channel, and hardens to form a solid skeletal structure. A wall can be formed of several block and channel levels. Kinard does not provide a strong external surface suitable for an outside wall.\nDunker, U.S. Pat. No. 4,624,089, issued on Nov. 25, 1986, is an anchor in the form of a plate having holes through which reinforcing rods extend, for holding reinforced sandwich panels together. An end of the anchor also wraps around a perpendicular mesh rod. Concrete is poured to form a carrier panel such that ends of several such anchors protrude therefrom. Then a layer of insulating material is fitted against the panel so that the protruding ends of the anchors punch through the insulating layer. Then rods for the mesh are fit through holes in the anchor protruding ends and another concrete layer is poured around the mesh and the anchor ends. This anchor holds two concrete layers around an insulating layer. Only the outside wall of most buildings needs to be this strong, so Dunker is wasteful of materials and needlessly expensive.\nMarks, U.S. Pat. No. 4,974,381, issued on Dec. 4, 1990, describes another anchor member for joining the outer slabs of a sandwich panel. A metal slat has a hole in each end for receiving a reinforcing rod of each slab, and the slat extends through a middle insulating layer. This anchor is not intended to anchor the panel to a foundation, but is instead part of the internal panel structure. Marks, like Dunker, is another wasteful sandwich arrangement.\nPorter, U.S. Pat. No. 4,947,600, issued on Aug. 14, 1990, teaches an interface for mounting a brick wall covering on an existing concrete wall or slab. One side of an angled member is attached to slab to form a shelf, additionally secured by a bracket. Studs extend through and join a foam layer to the slab. Mesh is placed over the foam layer and brick is laid on the shelf in front of the foam. Porter does not teach an effective, economical approach to constructing a new, insulated wall.\nHuettemann, U.S. Pat. No. 4,841,702, issued on Jun. 27, 1989, discloses a three-layer panel. The middle panel is an insulating slab such as foam polystyrene. A sheet of particle board is joined to one side and grooves are cut into the other side of the foam slab. Concrete is poured over the grooved side so that the concrete fills the grooves and creates structural ribs. Pouring continues until a layer of concrete is formed on top of the foam slab. Reinforcing rods may be placed in the grooves to strengthen the ribs. No efficient anchoring or lateral connection means are provided.\nZimmerman, U.S. Pat. No. 4,751,803, issued on Jun. 21, 1988, describes a multi-layer insulating panel having preformed concrete ribs. The ribs, referred to as studs, have metal connecting members protruding from one edge. The ribs are placed in parallel relationship in a jig with the connecting members protruding upward. Additional ribs are formed to create a border around the inside of the jig. A layer of insulating material is placed on top of the ribs and the protruding members pierce through the insulation material. Wire mesh is laid over the insulating material and concrete poured on top of the mesh. The concrete hardens around the connecting members but does not form a true unitary structure.\nOther references include Brostrom, U.S. Pat. No. 1,773,168 issued on Aug. 19, 1930; Branson, U.S. Pat. No. 1,840,304 issued on Jan. 12, 1932; Henzel, U.S. Pat. No. 2,321,813 issued on Jun. 15, 1943; Carroll, U.S. Pat. No. 4,090,336 issued on May 23, 1978; Lee, U.S. Pat. No. 4,380,887 issued on Apr. 26, 1983; Walston, U.S. Pat. No. 4,512,126 issued on Apr. 23, 1985; Pardo, U.S. Pat. No. 4,815,243, issued on Mar. 28, 1989; and Linetskii, Russian Patent Number 514,941 issued on Jun. 3, 1976.\nIt is thus an object of the present invention to provide a prefabricated panel having a simple design which is inexpensive to construct.\nIt is another object of the present invention to provide such a panel which has an exterior face of high strength and has superior insulating properties.\nIt is still another object of the present invention to provide such a panel which is simple to pour and permits rapid fabrication.\nIt is finally an object of the present invention to provide such a panel which has easy anchoring and interlock means for rapid, strong assembly."}
{"text": "Loading and unloading of road trailer trucks transporting containers is an important part of the operations in nearly all container terminals. As a function, the unloading of trailer trucks does not differ much from a situation where the container is picked up in the container storage area, as in both cases an unladen spreader is lowered on top of the container, after which the container is locked to the spreader by means of twistlocks.\nIn automatic container crane systems, the container can be automatically lifted from a trailer truck for example by means of measuring systems based on laser scanners. Systems have also been developed for loading the containers on a trailer automatically. An absolutely essential requirement for these systems is, however, that the desired position of the container can be accurately measured. However, accurate measurement requires reference points that can be identified. In most terminals, unloading and/or loading of trailer trucks is, nevertheless, a manual procedure due to security regulations.\nThe loading of trailer trucks is considerably more challenging than the unloading. The trailer in road trailer trucks transporting containers is usually provided with twistlocks, by means of which it is ensured that the container stays in place during the transportation. Lowering the container with sufficient accuracy in such a way as to align corner castings of the lowered container sufficiently accurately with the twistlocks requires accuracy. The visibility of the crane operator is disturbed by the container hanging from the spreader and blocking the twistlocks in the truck at the final stage of the lowering process.\nHandling of the trailers is facilitated by a variety of systems. For example, there are systems where the driver of a truck is assisted in stopping relative to the crane by means of a measuring system based on a laser scanner in such a way that the main movements of the crane would not have to be operated. By the main movements is meant the movement of the entire gantry or trolley.\nAn example of state of the art solutions representing laser scanner technology is patent publication CN201161875Y disclosing a system for determining locations of a spreader in a container crane and of a trailer truck to be loaded or unloaded, and particularly of twistlocks on a bed thereof (or of corner castings in a container) relative to the crane. Information produced by laser scanners in the system also allows one to conclude the type of a truck to be loaded and/or of a container to be lifted out from a truck.\nAnother example of solutions utilizing the laser scanner technology is disclosed in patent publication CN1884034A relating to a system for determining locations of a trailer truck loaded or unloaded by a container crane and of the containers on a bed thereof relative to the crane. Information produced by laser scanners in the system also allows one to conclude the type of a truck to be loaded and/or of a container to be lifted out from a truck.\nA third example of solutions representing this technology is patent publication JP2005239343A relating to a system for determining locations of trailer trucks to be loaded or unloaded relative to a crane handling containers. The system uses laser scanners disposed on the frame of the crane.\nA container can be automatically lifted from a trailer truck for example by means of the measuring systems based on laser scanners.\nThere are also systems based on cameras, wherein the aim is to recognize the desired position of the container of twistlocks in the trailer utilizing machine vision. One such solution is disclosed in patent publication DE3606363A1 relating to a system for determining the location of a trailer truck to be loaded or unloaded relative to a crane handling containers. The system uses cameras disposed on the frame of the crane at a height of a bed in the trailer truck. A problem of the camera-based machine vision systems is moderate reliability due to environmental conditions. In particular, wrong interpretations are caused by lighting that changes. When directing a spreader of a crane based on information from a measuring system, reliability is extremely important. False measurement data may cause material damages.\nYet another example of known state of the art is patent publication EP1337454A1, wherein a system for determining the location of a trailer truck to be loaded or unloaded relative to a crane handling containers is disclosed. The system uses sensors disposed on the frame of the crane at a height of a bed in the trailer truck.\nHowever, not all trailers have measurable reference points, such as twistlocks; instead, cargo is fastened to the trailer e.g. by tie down straps or chains. In addition, the twistlocks in some trailers are not exposed until points for the fastening of a container are subjected to the weight of the container."}
{"text": "For the communication of data packets, packet radio communications systems such as the General Packet Radio Service (GPRS) haven been developed. GPRS provides support for a packet-orientated services and is arranged to optimise network and radio resources for packet data communications. For example, GPRS networks can provide a facility for supporting internet protocol services to mobile user equipment. The GPRS provides a logical architecture, which is related to the circuit switched architecture of a mobile radio system.\nInternet protocol communications have become prevalent as a means for communicating data efficiently and conveniently. However the growth of internet protocol based services has introduced new demands on networks which operate in accordance with GPRS. It is therefore desirable to enhance packet radio communications systems, like GPRS, to make them more flexible and better able to cope with a rapid growth in internet protocol data traffic and internet protocol based services."}
{"text": "1. Field of the Invention\nThe invention relates to a keyboard assembly, more particularly to an ergonomical keyboard assembly.\n2. Description of the Related Art\nReferring to FIG. 1, a conventional keyboard assembly includes a plurality of keys which are arranged on a same plane, and an overall key surface formed on top of the keys is planar. When operating the keys of such an assembly, the user must place his/her hands in a working position in which palms and fingers lie on a substantially horizontal plane. In order to reach this working position, the user must turn or twist outward his/her elbows and wrists from their natural physical positions, and the palms, which normally face inward and towards one another and which are normally in a non-stretched convex state, have to turn to a downward facing position. Therefore, the aforesaid keyboard assembly is not ergonomically favorable to the user, and the user would easily become tired or even get hurt, especially, at his/her shoulders, elbows, and wrists, after working for a long period in an unnatural physical position.\nThere is another conventional keyboard assembly in which operating keys are divided into two groups (a, b) arranged in inclined positions in two separate areas, as shown in FIGS. 2 and 3, and in which an arched palm rest portion (c) is provided to support the hands of the user. This assembly permits the user to place his/her hands comfortably on the arched palm rest portion (c) in a natural physical position, without turning or twisting his/her elbows and wrists. However, since the key groups (a, b) are fabricated and are assembled separately, the manufacturing process of this keyboard assembly requires extra components and steps, thereby increasing the cost of production. On the other hand, as the overall key surface of each key group (a or b) is still planar and not arched, all of the keys must be hit in directions perpendicular to the planar key surface. As the hands have the fingers and palms thereof normally bending inward and as the fingers in their bending positions would have to move forward along directions perpendicular to an arched surface, this keyboard assembly is insufficient to satisfy the requirements for proper human operation and performance."}
{"text": "This invention relates generally to data communications, and, more particularly to a method and system for real-time reassembly of Asynchronous Transfer Mode (ATM) data.\nIn general terms, a computer network is a collection of end systems (also known as nodes) interconnected through one or more communication links. Generally, the end systems both send data (messages) to other end systems on the network and receive data sent by other end systems on the network.\nAmong available digital switching technologies, Asynchronous transfer mode (ATM) provides a promise of greater integration of services and increased flexible utilization of the network. Asynchronous transfer mode (ATM) configured networks allow high-speed data, voice and video communications to be conducted between endpoint computer systems. ATM networks, which are based on the transmission of fixed-length data packets, have proven to be extremely useful because they combine the benefits of both a switched network (i.e., constant transmission delay, guaranteed capacity) and a packet switched network (i.e., flexibility and efficiency for intermittent traffic).\nIn applications utilizing ATM, messages at a source are first broken up into these fixed length packets (ATM cells), transmitted, and then reassembled at a destination. ATM cells are 53 bytes long. They consist of a 5-byte header (containing an identifier of data flow which implicitly identifies the source address and the destination address) and a 48-byte information field. The header of an ATM cell contains all the information the network needs to relay the cell from one node to the next over a pre-established route. User data is contained in the remaining 48 bytes.\nATM uses a concept of virtual networking (or channels) to pass traffic between two locations, establishing virtual connections between a pair of ATM end-systems which are needed to connect a source with a destination. These connections are termed “virtual” to distinguish them from dedicated circuits. For any message, ATM cells always traverse the same path from source to destination. However, ATM does not have to reserve the path for one user exclusively. Any time a given user is not occupying a link, another user is free to use it.\nATM connections exist only as sets of routing tables held in each network node, switch, or other intermediate system, based on the virtual circuit identifier (VCI) and virtual path identifier (VPI) contained in the cell header. When a virtual path is established, each node (or switch) is provided with a set of lookup tables that identify an incoming cell by header address, route it through the node to the proper output port, and overwrite the incoming VCI/VPI with a new one that the next node along the route will recognize as an entry in its routing table.\nThe cell is thus passed from switch to switch over a prescribed route, but the route is “virtual” since the facility carrying the cell is dedicated to it only while the cell traverses it. Two cells that are ultimately headed for different destinations may be carried, one after the other, over the same physical wire for a common portion of their journey.\nTypically, an ATM Network Interface Card (NIC) and accompanying software are provided within the sending (or receiving) endpoint computer systems to transmit (or receive) the cells of a message over a virtual circuit. On the receiving end of the ATM network is another endpoint computer. This endpoint computer also includes a NIC that receives the incoming cells in the same order as they were transmitted. As the cells are received, it is the task of NIC and/or processor to reassemble the cells into the original message.\nConventional NICs can be described according to the reassembly method utilized. In one reassembly method, the NIC provides receives cells, places the received cells in a FIFO buffer, processes the cell information, and provides the processed cell information to the system memory. The NIC manages reassembly buffers in system memory, sending a received cell to the appropriate buffer based on its virtual circuit identifier and thereby reassembling the message. However, it is the responsibility of the user (and the device driver) to provide sufficient system bus bandwidth so that a substantially fixed rate of transfer (between the NIC and the system memory) is maintained. This type of reassembly is used, for example, in the IDT77252 product available from Integrated Device Technology Inc. of Santa Clara, Calif. Another reassembly method attempts to reduce the burden on the interconnecting devices and buses by providing enough memory on-board the NIC to store a complete message before transfer. Thus, with this type of NIC the entire message is reassembled in the NIC's local memory and then is transferred at one time from the NIC to the host memory.\nIn the reassembly techniques described in U.S. Pat. No. 6,097,734, reassembled cells are collected in groups of variable (programmable) size and sent to the host computer.\nIn the above described techniques, the reassembly does not occur in real-time. In order to monitor the performance of the network, an instrument utilizing real-time reassembly is desirable since such an instrument would introduce minimal delays in the transmission through the network. There is a need for methods and systems for reassembling, in real-time, the data received over an ATM network.\nIt is therefore an object of this invention to provide methods and systems for the real-time reassembly of data received over an ATM network.\nIt is a further object of this invention to provide systems for the real-time reassembling of messages from cells received over multiple ATM links."}
{"text": "Magnetic-drive centrifugal pumps are well suited for pumping caustic and hazardous fluids because shaft seals are not required. Instead of shaft seals, magnetic-drive pumps generally feature a pump shaft separated from a drive shaft by a containment shell. The drive shaft is arranged to rotate with a first magnetic assembly, which is magnetically coupled to a second magnetic assembly. The second magnetic assembly applies torque to the pump shaft to pump a fluid contained within the containment shell.\nThe reliability of containment shells may be rated in terms a burst strength. The burst strength is a pressure per unit area from the pumped fluid on the containment shell that results in damage to the containment shell sufficient to cause the leakage of fluid from the containment shell. In general, the higher burst strength, the better the containment shell. However, increasing the burst strength of a containment shell poses some difficult technical obstacles. For example, increasing the thickness of the containment shell or adding metallic reinforcement to the containment shell may significantly degrade pump performance, making any increase in the strength of the containment shell irrelevant. If the thickness of the containment shell is too great or if metal reinforcements are used indiscriminately, magnetic coupling between the first magnetic assembly and the second magnetic assembly may be impaired. In turn, the impeller may stop rotating entirely or may rotate too slowly for proper pump performance. Thus, a need exists for a containment shell with a superior burst strength, without sacrificing the requisite efficiency of the magnetic coupling between the first magnetic assembly and the second magnetic assembly.\nMany magnetically driven pumps include a front support and a rear support to support a rotating or a stationary pump shaft. The front support is often located such that the front support obstructs the inlet flow to the impeller, detrimentally limiting the performance of the pump under conditions of low net positive suction head (NPSH). Meanwhile, the rear support may be integral with a containment shell of polymer composite construction. The containment shell of the dual-support pump is often structurally inadequate to support a shaft without the assistance of a front support. Consequently, elimination of a front support for low net NPSH applications may reduce the burst strength of the containment shell, provide inadequate radial support for the pump shaft, or otherwise detrimentally impact pump reliability. Thus, a need exists for a containment shell that can support radial loads from a cantilevered shaft, while meeting a burst strength design goal."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the field of corn breeding. In particular, the invention relates to corn seed and plants of the variety designated CV338423, and derivatives and tissue cultures thereof.\n2. Description of Related Art\nThe goal of field crop breeding is to combine various desirable traits in a single variety/hybrid. Such desirable traits include enhanced yield, improved stalks, improved roots, resistance to insecticides, herbicides, pests, and disease, tolerance to heat and drought, reduced time to crop maturity, better agronomic quality, higher nutritional value, and uniformity in germination times, stand establishment, growth rate, maturity, and fruit size.\nBreeding techniques take advantage of a plant's method of pollination. There are two general methods of pollination: a plant self-pollinates if pollen from one flower is transferred to the same or another flower of the same plant. A plant cross-pollinates if pollen is transferred to it from a flower on a different plant.\nCorn plants (Zea mays L.) can be bred by both self-pollination and cross-pollination. Both types of pollination involve the corn plant's flowers. Corn has separate male and female flowers on the same plant, located on the tassel and the ear, respectively. Natural pollination occurs in corn when wind blows pollen from the tassels to the silks that protrude from the tops of the ear shoot.\nPlants that have been self-pollinated and selected for type over many generations become homozygous at almost all gene loci and produce a uniform population of true breeding progeny, a homozygous plant. A cross between two such homozygous plants produces a uniform population of hybrid plants that are heterozygous for many gene loci. Conversely, a cross of two plants, each heterozygous at a number of loci, produces a population of hybrid plants that differ genetically and are not uniform. The resulting non-uniformity makes performance unpredictable.\nThe development of uniform corn plant hybrids requires the development of homozygous inbred plants, the crossing of these inbred plants, and the evaluation of the crosses. Pedigree breeding and recurrent selection are examples of breeding methods used to develop inbred plants from breeding populations. Those breeding methods combine the genetic backgrounds from two or more inbred plants or various other broad-based sources into breeding pools from which new inbred plants are developed by selfing and selection of desired phenotypes. The new inbreds are crossed with other inbred plants and the hybrids from these crosses are evaluated to determine which of those have commercial potential.\nNorth American farmers plant tens of millions of acres of corn at the present time and there are extensive national and international commercial corn breeding programs. A continuing goal of these corn breeding programs is to develop corn hybrids that are based on stable inbred plants that have one or more desirable characteristics. To accomplish this goal, the corn breeder must select and develop superior inbred parental plants."}
{"text": "1. Field\nExample embodiments relate to a particle processing device. More particularly, example embodiments relate to a particle processing device using centrifugal force.\n2. Description of the Related Art\nGenerally, in order to process a fluid having particles, a particle processing device may include various pressure sources such as a syringe pump, a peristaltic pump, a gas pump, etc. However, an expensive external precise pump or a complex integrated micro-pump may be required for these types of the pressure sources.\nRecently, technologies for transferring a fluid by using a centrifugal force of the rotation of a substrate including a fluidic path formed therein have been developed. The particle processing device using the technologies may be manufactured by simple processes without the expensive micro-pump, to thereby reduce the manufacturing cost. However, since the conventional particle processing device is dedicated only for the fluid transfer, an actuator or valve using a pneumatice pressure or surface tension may be required for the device to control the fluid flow in order to process particles in the fluid."}
{"text": "The bounce angle of a golf club can impact the flight distance and/or accuracy of the golf ball, but the bounce height (i.e., the rise in the bounce angle) can have an even greater impact on the flight distance and/or accuracy of the golf ball. Specifically, the bounce angle and/or bounce height can determine how easily a golf club head can penetrate the ground under the golf ball during a golf club swing. However, ground conditions are subject to change such that the ground may be harder or softer at certain times due to weather, grooming, etc. Accordingly, a club head with a particular bounce angle and/or bounce height configuration may be better in certain circumstances than in other circumstances.\nFor simplicity and clarity of illustration, the drawing figures illustrate the general manner of construction, and descriptions and details of well-known features and techniques may be omitted to avoid unnecessarily obscuring the invention. Additionally, elements in the drawing figures are not necessarily drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help improve understanding of embodiments of the present invention. The same reference numerals in different figures denote the same elements.\nThe terms “first,” “second,” “third,” “fourth,” and the like in the description and in the claims, if any, are used for distinguishing between similar elements and not necessarily for describing a particular sequential or chronological order. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments described herein are, for example, capable of operation in sequences other than those illustrated or otherwise described herein. Furthermore, the terms “include,” and “have,” and any variations thereof, are intended to cover a non-exclusive inclusion, such that a process, method, system, article, device, or apparatus that comprises a list of elements is not necessarily limited to those elements, but may include other elements not expressly listed or inherent to such process, method, system, article, device, or apparatus.\nThe terms “left,” “right,” “front,” “back,” “top,” “bottom,” “over,” “under,” and the like in the description and in the claims, if any, are used for descriptive purposes and not necessarily for describing permanent relative positions. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the invention described herein are, for example, capable of operation in other orientations than those illustrated or otherwise described herein."}
{"text": "The present invention relates to an ink jet head, and to a method for the manufacture thereof.\nInk jet recording devices making utilization of the piezoelectric effect of piezoelectric elements to effect printing have been used as a printer, a word processing machine, a facsimile machine, etc. Some ink jet recording devices employ an ink jet head having a multi-level layered construction of thin films, like the one shown in FIGS. 14 and 15.\nReferring to the ink jet head-of the figures, formed in the interior of a main body portion 108 of the head are a plurality of pressure chambers 101 which are arranged along a direction L2, and a common ink chamber 102 which has a long and narrow shape in the direction L2. Each of the plurality of pressure chambers 101 is formed into a strip-like shape having a horizontal cross section which has a long and narrow shape in a direction L1. The pressure chambers 101 are in communication through their respective ink supply passages 105 with the common ink chamber 102. Nozzle orifices 103 are formed in a front side of the ink jet head (the lower one of the sides of the ink jet head in FIG. 14). The nozzle orifice 103 is in communication through an ink discharging passage 104 with the pressure chamber 101. Each pressure chamber 101 is compartmented with an oscillation plate 112 at its back side (the upper one of the sides in FIG. 14). The oscillation plate 112 is formed of Cr (chromium) and functions as a common electrode for applying a voltage to piezoelectric elements 106. The piezoelectric element 106 is jointed to a back side of the oscillation plate 112. Mounted on a back side of the piezoelectric element 106 is an individual electrode 111. The individual electrode 111, the piezoelectric element 106 and the oscillation plate 112 together form an actuator portion 107 for applying a pressure to the ink in the pressure chamber 101. The main body portion 108 and the oscillation plate 112 are bonded together by an adhesive 120.\nWith the recent advancement of the resolution of recording devices to a higher degree, the miniaturization of head structures has been improved. Correspondingly, the gap between two piezoelectric elements 106 next to each other becomes shorter as shown in FIG. 15 and the thickness of a compartmenting wall 109 which compartments the pressure chamber 101 becomes thinner. Accordingly, when adhering together the main body portion 108 and the oscillation plate 112, there are some cases in which the adhesive 120 is forced out of the butt surface, as indicated by the virtual line (the two-dot chain line) of FIG. 15 and a part 114 of the forced-out adhesive adheres to a displacement portion 113 of the oscillation plate 112. This causes the displacement portion 113 of the oscillation plate 112 to decrease in its displacement amount, in addition to which the displacement portion 113 is stopped from performing a smooth displacement operation, therefore resulting in a drop in the ink discharging performance. Further, it is possible that the displacement amount of the oscillation plate 112 varies from one pressure chamber to another, resulting in poor print quality.\nBearing in mind the above-described problems with the prior art techniques, the invention was made. Accordingly, an object of the present invention is to provide such adhesion between a head main body and an oscillation plate that the adhesive will not become an obstacle to the displacement operation of the oscillation plate.\nIn order to achieve the object, the present invention provides an ink jet head. More specifically, the ink jet head of the present invention comprises (a) a head main body portion in which a pressure-chamber recess portion and a nozzle orifice are formed and (b) an actuator portion having (i) an oscillation plate which covers the pressure-chamber recess portion of the head main body portion so as to form, together with the head main body portion, a pressure chamber in the form of a compartment and (ii) a piezoelectric element which is rigidly fixed to the oscillation plate whereby the displacement of the oscillation plate generated by the piezoelectric effect of the piezoelectric element forces ink in the pressure chamber out of the nozzle orifice, wherein an intermediate layer is formed between the oscillation plate and the head main body portion for the prevention of adhesive adhesion to the oscillation plate and wherein the oscillation plate and the head main body portion are rigidly fixed together by an adhesive through the intermediate layer.\nAs a result of such arrangement, even when a part of the adhesive is forced out in adhering together the head main body and the oscillation plate, the provision of the intermediate layer between the adhesive and the oscillation plate causes the forced-out adhesive to adhere to the intermediate layer. Accordingly, the forced-out adhesive is prevented from adhering to the displacement portion of the oscillation plate, therefore ensuring that the oscillation plate satisfactorily performs its displacement operation. As a result, the variation in displacement is prevented, thereby achieving a stable ink discharging performance.\nThe present invention provides another ink jet head which comprises (a) a head main body portion in which a pressure-chamber recess portion and a nozzle orifice are formed and (b) an actuator portion having (i) an oscillation plate which covers the pressure-chamber recess portion of the head main body portion so as to form, together with the head main body portion, a pressure chamber in the form of a compartment and (ii) a piezoelectric element which is rigidly fixed to the oscillation plate whereby the displacement of the oscillation plate generated by the piezoelectric effect of the piezoelectric element forces ink in the pressure chamber out of the nozzle orifice, wherein an intermediate layer having a window portion through which a displacement portion of the oscillation plate is exposed to the pressure chamber, and consisting of metal, ceramics, or resin is formed on a surface of the oscillation plate on the side of the head main body, and wherein the oscillation plate and the head main body portion are rigidly fixed together by an adhesive through the intermediate layer.\nAccordingly, by the use of thin film formation technology including sputtering and CVD, the intermediate layer used for the prevention of adhesive adhesion to the displacement portion of the oscillation plate can be formed easily.\nThe use of copper facilitates the formation of relatively thick films, and therefore the intermediate layer is preferably formed of copper. As a result of such arrangement, it is possible to easily form the intermediate layer having a thicker film thickness capable of the prevention of adhesive adhesion to the oscillation plate without fail. Moreover, the copper is superior in workability, that is, it is easily etched. Therefore, the window portion can be formed easily.\nThe intermediate layer may be formed of titanium. This extends the life span of the intermediate layer because the titanium is not readily degraded even when exposed to the ink. Moreover, the degree of adhesion between titanium and chromium is high, so that, when chromium is used for the oscillation plate, the adhesion with respect to the oscillation plate is enhanced. Accordingly, there is no need for the provision of an adhesion improving layer used to improve the degree of adhesion between the intermediate layer and the oscillation plate, whereby the manufacture process can be shortened.\nIf the intermediate layer is formed too thin, then the adhesive is likely to adhere to the displacement portion of the oscillation plate. Taking into account this point, it is preferred that the intermediate layer is formed, having a thickness of 5 xcexcm or greater. As a result of such arrangement, it becomes possible to effectively prevent the adhesive from adhering to the displacement portion of the oscillation plate.\nFurthermore, it is more preferable that the intermediate layer is formed, having a thickness of more than 7 xcexcm. This further ensures that the adhesive is prevented from adhering to the displacement portion of the oscillation plate.\nOn the other hand, if the intermediate layer is formed too thick, this makes it difficult to form the intermediate film having a well arranged shape. It is therefore preferable that the intermediate layer is formed, having a thickness of 30 xcexcm or less. This facilitates the formation of intermediate layers having an adequate shape.\nThe degree of adhesion between the intermediate layer and the oscillation plate become insufficient in some cases because of the material used. Therefore, an arrangement may be made in which an adhesion improving layer is sandwiched between the intermediate layer and the oscillation plate for the enhancement of the degree of adhesion between them. By virtue of the placement of the adhesion improving layer between the intermediate layer and the oscillation plate, the degree of adhesion between them can be enhanced, as a result of which the peeling-off of the intermediate layer can be prevented.\nIt is arranged such that the oscillation plate is formed of chromium, the intermediate layer is formed of copper, and the adhesion improving layer is formed of titanium. Because of such arrangement, even though the oscillation plate is formed of chromium and the intermediate layer is formed of copper, it is possible to tightly adhere the intermediate layer and the oscillation plate together because the degree of adhesion of titanium with respect to chromium is high.\nIn order to prevent the adhesion improving layer from being formed into an island-like shape so that the adhesion improving layer becomes a satisfactory layer, it is preferred that the adhesion improving layer is formed, having a thickness of 0.01 xcexcm or greater. As a result of such arrangement, the adhesion improving layer functions to tightly adhere the oscillation plate and the intermediate plate together.\nThe present invention provides still another ink jet head which comprises (a) a head main body portion in which a pressure-chamber recess portion and a nozzle orifice are formed and (b) an actuator portion having (i) an oscillation plate which covers the pressure-chamber recess portion of the head main body portion so as to form, together with the head main body portion, a pressure chamber in the form of a compartment and (ii) a piezoelectric element which is rigidly fixed to the oscillation plate whereby the displacement of the oscillation plate generated by the piezoelectric effect of the piezoelectric element forces ink in the pressure chamber out of the nozzle orifice, wherein a projection portion for mounting the oscillation plate on the head main body portion is formed on a surface of the oscillation plate opposite to the head main body portion.\nAn arrangement may be made, in which the oscillation plate and the head main body portion are rigidly fixed together by an adhesive through the projection portion.\nAccordingly, the projection portion which is a part of the oscillation plate functions as the foregoing intermediate layer. As a result of such arrangement, the intermediate layer is no longer to be formed separately from the oscillation plate. Moreover, there is no need to attach the intermediate layer to the oscillation plate.\nThe present invention provides a method of manufacturing an ink jet head in which an oscillation plate is displaced by the piezoelectric effect of a piezoelectric element and the displacement of the oscillation plate forces ink in a pressure chamber out of a nozzle orifice, the method comprising the steps of (a) forming on a surface of the oscillation plate an intermediate layer, (b) forming an opening in a position of said intermediate layer corresponding to a displacement portion of the oscillation plate, and (c) with the opening of said intermediate layer brought into agreement in position with a pressure-chamber recess portion of the head main body portion, adhering together the intermediate layer and the head main body portion.\nAs a result of such arrangement, the oscillation plate and the head main body portion are rigidly adhered together by adhesive through the intermediate layer. Accordingly, even when a part of the adhesive is forced out in adhering the oscillation plate and the head main body portion together, the forced-out adhesive will adhere to the intermediate layer, whereby adhesive adhesion to the displacement portion of the oscillation plate is prevented.\nThe present invention provides another method of manufacturing an ink jet head in which an oscillation plate is displaced by the piezoelectric effect of a piezoelectric element and the displacement of the oscillation plate forces ink in a pressure chamber out of a nozzle orifice, the method comprising the steps of (a) forming on a surface of the oscillation plate an adhesion improving layer, (b) forming on a surface of the adhesion improving layer an intermediate layer, (c) forming an opening in at least a position of the intermediate layer corresponding to a displacement portion of the oscillation plate, and (d) with the opening of the intermediate layer brought into agreement in position with a pressure-chamber recess portion of the head main body portion, adhering together the intermediate layer and the head main body portion.\nAs a result of such arrangement, the adhesion improving layer is formed between the oscillation plate and the intermediate layer, wherein the intermediate layer is rigidly secured to the oscillation plate through the adhesion improving layer. Thus, the oscillation plate and the head main body portion are tightly secured together by an adhesive through the intermediate layer. Accordingly, even when a part of the adhesive is forced out in adhering the oscillation plate and the head main body portion together, the forced-out adhesive will adhere to the intermediate layer, whereby adhesive adhesion to the displacement portion of the oscillation plate is prevented.\nAs described above, in accordance with the present invention, the intermediate layer is formed between the oscillation plate and the head main body portion, so that the adhesive is prevented from adhering to the displacement portion of the oscillation plate in adhering the actuator portion and the head main body portion together. This accordingly enables the oscillation plate to perform a displacement operation in a smooth manner, thereby achieving an improved ink discharging performance.\nThe intermediate layer is formed of copper, thereby facilitating the formation of a relatively thick film well-arranged in shape. This therefore makes it possible to form, in an easy manner, a layer suitably used for the prevention of adhesive adhesion to the oscillation plate. Moreover, it is possible to easily form the window portion for the displacement portion of the oscillation plate.\nThe arrangement that the intermediate layer is formed of titanium makes it possible to make the intermediate layer less readily degraded when exposed to ink. Moreover, if the oscillation plate is formed of chromium, this makes it possible to enhance the degree of adhesion between the intermediate layer and the oscillation plate.\nBy virtue of the provision of the adhesion improving layer between the intermediate layer and the oscillation plate, the degree of adhesion between the intermediate layer and the oscillation plate can be enhanced, therefore preventing the peeling-off of the intermediate layer without fail."}
{"text": "A refrigeration system in a retail store, such as, for instance, a supermarket or grocery store, may include a large number of refrigeration circuits, such as, for instance, refrigeration and freezer display cases and walk-in coolers. A refrigeration system may also include one or more rooms, commonly referred to as refrigeration racks, which may contain compressors, fans, refrigerant tanks, and/or associated control circuitry, to centrally manage the temperature of the refrigeration circuits (e.g., keep the refrigeration circuits at the correct temperature). For example, the refrigeration system may use refrigerant held in the refrigerant tanks to transfer heat from one location to another in the refrigeration system.\nA refrigerant leak may occur in a refrigeration system at some point during operation of the system. If the refrigerant leak is not quickly detected, an unplanned loss of refrigerant will occur, which may result in a refrigeration rack failure and/or emergency situation in the retail store."}
{"text": "Various prior art structures and organizations have been employed relative to transport of cargo within a pickup truck bed, wherein such is exemplified by U.S. Pat. Nos. 5,110,172 and 5,149,164 relative to pickup truck bed tailgates and liner structure.\nThe U.S. Pat. No. 3,961,585 is directed to a cargo net that is arranged for surrounding and securing cargo relative to a pallet, with the U.S. Pat. No. 5,275,458 indicating construction relative to a tonneau cover.\nU.S. Pat. No. 3,478,394 relative to a web hardware adjuster is provided and incorporated herein by reference as an example of strap adjusting structure."}
{"text": "The present invention relates to a loudspeaker system and especially to a loudspeaker system using an elongated hollow tube having wide range speakers mounted in each end thereof for emanating sounds from each of the tube.\nThe loudspeaker tube system in accordance with the present invention will produce a stereophonic sound or a dual monophonic sound system which is self-contained in a single length of pipe or tubing with a full range loudspeaker at either end of the length of the pipe. The design simplifies the otherwise cumbersome loudspeaker components normally associated with higher quality sound systems by having a single module that does all of the work for the system.\nIn the past, stereophonic sound systems of high quality sound reproductions have been incorporated into two or more loudspeaker enclosures. These enclosures employ more than one driver inside the loudspeaker box and many employ a woofer, midrange, and a tweeter driver in at least two enclosures separated by distance between the two speakers to achieve a sound stage. Sometimes a subwoofer is used to handle bass frequencies below 100 Hz.\nThe present invention is a full range loudspeaker system which maximizes the bass frequencies in the subwoofer range and isobaric configurations set the two drivers out of phase with each other in the tube to increase the lower frequency sound pressure levels while minimizing box size requirements. This technique is commonly called \"push-pull\" because one loudspeaker pushes while the other pulls using the air pressure in the enclosure. This technique has been successfully used in low frequency units in the past but has had severe limitations in full range loudspeakers. One problem in full range loudspeakers has been the phase shift in a full range loudspeaker system, particularly where a stereophonic or discrete two channel signals are introduced. The phase must be correct or the sound stage will be off-center. In addition, one speaker will seem to have more bass than the other speaker, which will sometimes appear tinny or thin. An advantage of the present system using a length of tubing is that a discrete stereo signal can be placed in phase to each driver by having one driver face the furthest distance from the other inside a single pipe. The signal from channel one will cause the associated driver to push out thereby rarifying the air within the container and pulling the opposing driver along while doubling the SPL and bass frequency response. The second driver, which receives a second signal, will electrically be pulling in the opposite direction to that of the acoustic action caused by the first signal reaction. Thus, a dual opposing signal or extended stereo signal can create a push-pull affect on a full range loudspeaker system if that system is contained in a single unit where the drivers are set to the furthest point within the unit away from each other, in phase electrically, and where the enclosure is equal to a minimum of twice the length of the radius of one of the drivers and assuming both drivers are of equal character, size, shape, and general electrically acoustical response to one another.\nThere have also been small stereophonic systems in the past but the advantage of the cylindrical tube is that it is a singular and small pipe with a length of five inches or more, can put out more bass with more sound pressure level, and more full range frequency response with more clarity and with more perceived volume than systems three to four times its size and weight.\nIn the stereophonic full range mode, the sound stage can be easily adjusted using simple mechanical controls or walls to reflect the sound in the direction you would have it to go. The width of the sound stage can be spread out or brought in closer depending upon how far out you choose to place the simple deflection walls.\nAnother way to create a sound stage which incorporates an acoustic center field image is to cut a hole into the center of the pipe equal in size to at least half the radius of either of the two drivers and force a length of pipe, at least the length of the diameter of the hole, into the hole, closing off all air from escaping outside the center hole. The interaction between the two stereo signals using standard tuning procedures, like a helmholtz or a small theile tuned port, will produce the third product of sound which results from the tuned port or an acoustical center field.\nIn a dual mono-mode, where each of the two speakers are driven by a monophonic signal, a better bass signal can be produced and two pipes can be used driven by the same monophonic signal to produce a bigger, louder sound capability. In the dual monophonic configuration each driver can be wired in standard phase, which produces a single point source full range uniform response, or out-of-phase using the isobaric configuration. In this configuration, the bass response will increase below 150 Hz but the length of pipe needs to be cut to be equal to at least the length of one of the drivers radius. For instance, if the driver is 4\" around, the length between the two drivers would be approximately 4\" long. The disadvantages of this configuration is the sacrifice that some of the full range uniformity but where two separate pipes are employed, the bass extension or in-phase driver can be loaded towards the wall to increase the bass by 3 db per pipe with the full range sound extended into the room acting as the point source. The same system can be also accomplished with a longer pipe but the electroacoustic phasing will not be as effective for delivery of deeper bass as it would in the smaller enclosure using generally standard type drivers.\nPrior art speaker systems may be seen in the following U.S. patents. In the Ashe U.S. Pat. No. 2,905,259, a tubular speaker housing has a single loudspeaker mounted in the middle of two truncated cones attached to form one tube spreading in each direction. The Lanternier U.S. Pat. No. 4,905,788, shows an electro-acoustic transducer which uses a single hollow cylindrical support for an active speaker at one end and a passive radiator mounted at the opposite end thereof for increasing the bass response from the tube. Two systems are required for stereo response. In the Flanders U.S. Pat. No. 4,753,317, a trapezoidal loudspeaker enclosure mounts a loudspeaker at one end at an angle and passively at the other with attached base plates over each angled end of the enclosure. The Manger U.S. Pat. No. 4,268,719, is a loudspeaker arrangement giving a pair of back-to-back loudspeakers in an infinite baffle arrangement. The Seville U.S. Pat. No. 4,655,315, teaches a speaker system having a plurality of serially connected angularly disposed hollow tubes with a loudspeaker mounted at one end, which tube can be adjusted to different shapes. The Virva U.S. Pat. No. 4,164,988, shows a fine tuned column speaker system in which a bellows is formed in part of tube or in which telescoping tubes can be adjusted for a single loudspeaker to vary the length of the back tube for the backwave of the speaker. The Robinson U.S. Pat. No. 3,945,461, is a sound speaker system which has an adjustable telescoping tube which can be locked in place and has a loudspeaker mounted facing the bottom of one of the tubes.\nThe aim of the present invention is to have a loudspeaker system with a single elongated hollow tube which may be a rigid tube or a telescoping tube with a wide range speaker mounted in each end thereof and the tube length set to accomplish the desired affect in both the wide range sound and increased bass response and also to utilize positioned end deflectors for directing the sound."}
{"text": "It is necessary to input in advance into a silver halide photographic light-sensitive material various information related to the photographic light-sensitive material such as a kind and a production number of the photographic light-sensitive material, the name of manufacturer and the emulsion No., various information related to the photographing with a camera for example, date and time of photographing, aperture, exposure time, strobe conditions, filter used, climate conditions, size of photographing frame, a type of a camera and the use of anamolphic lens, various information necessary for printing, for example, the number of prints, choice of filter, color tasty of the customer and the size of trimming, various information obtained in printing, for example, the number of prints, choice of filter, color tasty of the customer and the size of trimming and other information related to customers from a viewpoint of management and from the viewpoint of an improvement of print quality and the enhancement of efficiency of printing operation.\nIn conventional photographic light-sensitive materials, it has been impossible to input all of the above information. Only information such as the date and the time of photographing, aperture and exposure time has been inputted optically. In addition, in printing, there has been no means for inputting the above-mentioned information into photographic light-sensitive material.\nIn a magnetic recording system, recording and playbacking are easy. Accordingly, magnetic recording systems for inputting the above-mentioned various information in photographic light-sensitive materials have been studied, and various technologies have been proposed.\nFor example, technologies to provide a stripe-type magnetic recording layer wherein fine particles of ferromagnetic substance are dispersed on an emulsion surface at the side of an image portion or a backing surface for recording information such as sound and photographing conditions are disclosed in Japanese Patent Publication Open to Public Inspection (hereinafter referred to as Japanese Patent O.P.I. Publication) Nos. 62627/1975 and 4503/1974 and U.S. Pat. Nos. 3,243,376 and 3,220,843, and, further, technologies to provide a transparent magnetic recording layer wherein necessary transparency is obtained by selecting the amount and the size of magnetic particles on the backing layer of photographic light-sensitive material are disclosed in U.S. Pat. Nos. 3,782,947, 4,279,945 and 4,302,523. In addition, a roll film having on its reverse side a magnetic recording layer containing a magnetic substance capable of recording magnetically and a photographing camera having a magnetic head as well are disclosed in U.S. Pat. No. 4,947,196 and WO90/04254.\nBy providing the above-mentioned magnetic recording layers, it has come to be possible to record in photographic light-sensitive materials the above-mentioned various information which had been difficult previously and it has come to be promising to record sound and image signals.\nOn the other hand, it has recently been demanded to reduce the size of a camera. Ordinarily, it is the most effective to reduce the size of a cartridge (magazine) further for the reduction of a camera size. For this, it is necessary to reduce the thickness of a support of a photographic light-sensitive material.\nHowever, when the thickness of a support is reduced, the mechanical strength of the support is lowered, and on some occasions, damages such as folding and breakage starting from a perforation portion are caused.\nSince a polyethylene terephthalate film is more excellent in terms of mechanical strength than a cellulose triacetate film, it is possible to reduce the thickness of the support.\nHowever, in the case that a polyethylene terephthalate film is used as a support of a photographic light-sensitive material having a magnetic recording layer, a roll set curl is caused when the film is stored in its shape of a roll at a high temperature. Since this roll set curl cannot be removed sufficiently in the course of photographic processing, contact with a magnetic head cannot be kept satisfactorily. Accordingly, there has been caused a problem of an occurrence of errors in the course of inputting magnetic signals and playbacking thereof.\nIn addition, when a magnetic recording layer is provided on a support, the magnetic recording layer of a photographic light-sensitive material comes, in contact with a magnetic head in the course of recording and playbacking information, causing the photographic light-sensitive material to be charged to cause an error in the course of playbacking the record.\nFurthermore, there has been a problem of adhesion of dust due to static, causing deterioration of a photographic print.\nFor preventing the above-mentioned problems, there have so far been known methods of providing a conductive layer containing a conductive material. Among these, the use of a crystalline metal oxide has been known useful because of little temperature dependence of conductivity as is disclosed in Japanese Patent O.P.I. Publication No. 143431/1981; however, on the reverse side for impeding the accumulation of static electricity as is disclosed in Japanese Patent O.P.I. Publication No. 62543/1992.\nIn this meaning, however, flying-off of static charge is not conducted sufficiently in a short time. Especially, when the speed of the conveyance of the photographic light-sensitive material is increased during recording and playbacking of information, an error occurs or the product value of photographic prints is lowered due to adherence of dust caused by static electricity.\nIn addition, when polyethylene terephthalate is used for assuring mechanical strength of a support, it is easily charged, making more serious problems to be caused occasionally."}
{"text": "1. Field of the Invention\nThe present invention relates to a hologram memory apparatus for recording and reconstructing spatial pattern information. More precisely, the present invention relates to an improvement of such an apparatus of which the structure is greatly simplified by using the one-dimensional hologram recording.\n2. Description of the Prior Art\nAs a conventional method and apparatus for recording information in a one-dimensional hologram, it is known to use a computer-generated hologram, in which an input information is Fourier transformed by means of a computer to produce an electric holographic signal and an electron beam or a laser beam is modulated by the electric signal and is used for recording. For instance, Adam Kozma et al had submitted a report in the 1971 IEEE/OSA Conference on Laser Engineering and Applications (Session No. 152). This suggested system has an advantage in that a spatial modulator is not required; however, it contains a disadvantage in that the writing speed cannot be made high enough owing to a long computation time.\nThere are various disclosures with respect to the two-dimensional page-oriented hologram memory. This kind of apparatus, however, had encountered a number of technical difficulties mainly due to the complicated construction of the structural devices since it requires various elements such as a two-dimensional spatial modulator, i.e. page composer or a two-dimensional photo detector matrix, etc.\nFurthermore, there had been an attempt to realize a simplified hologram memory system considering applications as wideband recorders. For instance, Andrew Bardos disclosed such a system in Applied Optics, Vol. 13, No. 4, April 1974, pp. 832-840. In this system, a time series of a given electrical pulse information is rearranged into multi-channel information and by using such multi-channel information the respective channel of a multi-channel acousto-optic light deflector is controlled as a spatial modulator and thus the input signal light beam is formed. The input signal light beam is subjected with a doppler shift. Accordingly, in the recording of the hologram in order to stabilize the interference fringe, the reference light should also be subjected with such a doppler shift to the same extent with the signal light by using another acousto-optic light deflector, names an acoustic beam splitter, to give diffractive deflection for the laser beam. Accordingly, the optical path for the reference light beam is separated from that of the signal light path at the acoustic beam splitter and thus the strip-shaped hologram recording is made through a writing lens after readjustment of the signal light beam and the reference light beam to be parallel with each other. There exist two problems in this step. One is the stability of the system due to use of a separate optical path. Another is the demand for large aperture lens systems, which make them costly and difficult to produce."}
{"text": "As the population of computer users grows, so does the variety of special needs of the users. For a user to maximize their efficiency in completing tasks on a computer, the user interface must meet their specific requirements.\nFor the purposes of this document, the term “typical computer user” will refer to individuals without any impairment, mild or severe. A typical computer user can effectively use a computer without employing any adaptive hardware or software, and without using any accessibility features that may be built into the computer's operating system.\nBlind computers users are unable to see the screen on a computer so must access the information using some form of access technology. They may listen to the information through a text to speech engine, read the information on a refreshable Braille display or feel the information utilizing a Haptic display. With these adaptive features, software delivers the information to the various output devices by taking the information from the screen buffer, acquiring it from an application program interface (API) or by generating the information internally from its own data.\nLow vision users differ from blind users as they can see some things but are still unable to read a computer screen without augmentation to the image being displayed. Low vision computer users may employ screen magnification software, screen reader software using text to speech output or work with the features of the operating system to change colors, mouse pointers and other attributes of the display to make customize the image to address their specific requirements.\nDeaf and hearing impaired computer users can see the contents of a computer screen but are unable to hear the information delivered in an auditory manner. Many deaf people also have difficulty with written text in their native language as the various sign languages used around the world often differ from the spoken or written word. To maximize their efficiency on a computer, a deaf or hearing impaired user may employ visual cues to deliver information that otherwise may have been presented as a sound. A deaf user may also use a screen reader that turns written text into sign language delivered by an animated avatar.\nComputer users who are both deaf and blind use computers with screen reading software that sends information to a refreshable Braille display. To interface with the computer, they provide input to the computer utilizing controls on the Braille display or with a keyboard.\nIndividuals with learning disabilities have a wide variety of different needs that must be met in order for them to effectively use a computer. These requirements vary from changing colors of text, changing the spacing between words, changing the shape of characters to be tall and thin or short and wide, using software that delivers information through a text to speech engine or employing a wide array of different augmentations to the information presented on a computer screen.\nMany mobility impaired users have no difficulty accessing the information provided by a computer but, rather, require physical accommodation, such as a special desk designed to accommodate a wheelchair. Individuals whose motor impairments affect their ability to type on a keyboard or otherwise use their hands do require special adaptive hardware and software to be able to provide input and command and control to a computer.\nFor the purposes of this document, the term, “temporarily disabled” refers to individuals who, due to a temporary circumstance, lose the ability to use their eyes, ears, hands or some other human attribute. Examples of temporarily disabled individuals include motorists, who cannot take their eyes off of the road or their hands off of the steering wheel; people working in very loud places where they cannot hear the audio from a computer and people performing other tasks that cause them to be unable to interact with a computer in a typical fashion.\nThe list above describes some of the more common special needs that a computer user may experience, either permanently or temporarily.\nVisual difficulties and impairments include low vision, color blindness, and blindness. There are many options for individuals with vision difficulties to modify the computer displays and appearance so it is more legible, or receive information through sound or touch. These options are commonly referred to in the industry as, “assistive” or “adaptive” technologies. Individuals who have visual difficulties and impairments may utilize any one, or a combination of, assistive technologies including screen enlargers, screen readers, speech recognition systems, voice recognition programs, speech synthesizers, text-to-speech systems, refreshable Braille displays, Braille embossers, talking word processors and large-print word processors.\nTo succeed in a professional environment, most employees need to create, read, collaborate on and share documents of many different types generated from many different software programs. Visually impaired individuals have often struggled to perform some of the most basic document related tasks as it is very difficult and in some cases, impossible for them to understand the subtleties conveyed through visual attributes and alignment of text.\nSystems currently known in the art utilize unstructured data formats to provide data to assistive technology devices. The use of unstructured data formats requires that the software make continuous calls to component object modules (COM) to establish the relationship between the assistive technology device and the unstructured data format. With these call to COM objects, the data formats are unable to convey the subtleties inherent in most basic documents. The need to make repeated calls to a COM object required substantial processing and generates latencies in the system that are effect the experience of the end user.\nIt is known in the art that computer programs utilize extensible markup language (XML) as the native document format for their data storage. XML is simply a string format in which the string contains text descriptions (descriptors) for the data in the string. These descriptors allow a parser to find information in the string without the string having a predefined format.\nAccordingly, what is needed in the art is a system and method that transforms XML information from an application into a user definable audio, visual or tactile enriched document."}
{"text": "The use of fossil fuel as the combustible fuel in engines results in the combustion products of carbon monoxide, carbon dioxide, water vapor, smoke and particulate, unburned hydrocarbons, nitrogen oxides and sulfur oxides. Of these above products carbon dioxide and water vapor are considered normal and unobjectionable. In most applications, governmental imposed regulations are restricting the amount of pollutants being emitted in the exhaust gases.\nIn the past, NOx emissions have been reduced by reducing the intake manifold temperature, retarding the injection timing, and modifying the injection rate shape. And, the adverse effects on fuel consumption, particulate emissions, and engine performance have largely been alleviated through improvements in the basic engine design and fuel selection. For example, at the present time smoke has normally been controlled by design modifications in the combustion chamber, particulates are normally controlled by traps and filters, and sulfur oxides are normally controlled by the selection of fuels being low in total sulfur. This leaves carbon monoxide, unburned hydrocarbons and nitrogen oxides as the emission constituents of primary concern in the exhaust gas being emitted from the engine.\nMany systems have been developed for recycling a portion of the exhaust gas through the engine thereby reducing the emission of these constituents into the atmosphere. The recirculation of a portion of exhaust gas is used to reduce pollution emitted to the atmosphere. In a naturally aspirated engine this process is relative simple. But, with a turbocharged engine, the recirculation of a portion of the exhaust gas into the intake air becomes more complex because the intake pressure is higher than the exhaust pressure during many operating conditions. In many of such past system a volume of the exhaust gas from the engine was redirected to the intake air of the engine through the turbocharger and/or an aftercooler and to the engine. Such systems caused the premature plugging of aftercooler cores and malfunctioning of the systems. Additionally, with such recirculation system deterioration of the exhaust flow was caused by deposit buildup.\nVarious approaches have been used to address the adverse pressure gradient issue. For example, throttling valves have been installed in the air inlet, back pressure valves in the exhaust gas, intake manifold venturis, etc. to provide sufficient pressure drop to get the exhaust gas to flow to the intake air. Although this provides the necessary pressure drop to functionally operate an exhaust gas recirculation system several disadvantages, such as, fuel consumption, emissions, and/or performance occur.\nThe present invention is directed to overcoming one or more of the problems as set forth above."}
{"text": "The present invention disclosed herein relates generally to memory systems, and more particularly, to memory systems and memory cards in which a bad block is used in a single level cell (SLC) mode.\nRecently, the number of devices using nonvolatile memories has increased. For example, an MP3 player, a digital camera, a mobile phone, a camcorder, a flash card, and a solid state disk (SSD) are examples of devices that use nonvolatile memories as storage devices.\nAs more devices use nonvolatile memories as storage devices, the capacity of the nonvolatile memory is drastically increasing. One method for increasing memory capacity is a so-called multi level cell (MLC) method in which a plurality of bits is stored in one memory cell.\nFIG. 1 is a block diagram illustrating a conventional memory system. Referring to FIG. 1, a conventional memory system 100 includes a host 110, a memory controller 120, and a flash memory 130.\nThe memory controller 120 includes a buffer memory 121. The flash memory 130 includes a cell array 131 and a page buffer 132. Although not shown in FIG. 1, the flash memory 130 includes a decoder, a data buffer, and a control unit.\nThe memory controller 120 receives data and a write command from the host 110, and controls the flash memory 130 to program data into the cell array 131. Also, the memory controller 120 controls the flash memory 130 to read data stored in the cell array 131 according to a read command input from the host 110.\nThe buffer memory 121 temporarily stores therein data to be programmed into the flash memory 130, and data read from the flash memory 130. The buffer memory 121 transfers the temporarily stored data to the host 110 or the flash memory 130 under control of the memory controller 120.\nThe cell array 131 of the flash memory 130 includes a plurality of cells. The memory cells are nonvolatile and can retain stored data even when no power is applied. A page buffer 132 is a buffer that stores data to be programmed into a selected page of the cell array, or data read from a selected page.\nA memory cell of the flash memory 130 is categorized into a single level cell (SLC) and a multi level cell (MLC) according to the number of data bits that can be stored therein. The SLC can store single-bit data, and the MLC can store multi-bit data.\nThe SLC, a memory cell storing single-bit data, will now be described. The SLC has two states depending on distribution of a threshold voltage. After being programmed, the memory cell stores data 1 or data 0. A memory cell storing data 1 is considered to be in an erase state, and storing data 0 is considered to be in a program state. The cell in the erase state may be called an “on-cell,” and the cell in the program state may be called an “off-cell.”\nThe flash memory 130 performs a program operation by page. The memory controller 120 transfers data to the flash memory 130 by page, using the internal buffer memory 121 during a program operation.\nThe page buffer 132 temporarily stores data loaded from the buffer memory 121, and simultaneously programs the loaded data to a selected page. After the program operation is completed, a program verification operation for verifying whether data is correctly programmed or not is performed.\nIf the program verification result is a “program fail,” a program operation and a program verification operation are performed again with an increased program voltage. After the program operation for one-page of data is completed in the above described manner, the next data are received for a program operation.\nAn MLC, a memory cell storing multi-bit data, will now be described. FIGS. 2 and 3 illustrate a process of programming a least significant bit (LSB) and a most significant bit (MSB), that is, 2-bit data, into one memory cell.\nReferring to FIG. 2, a memory cell is programmed to have any one of four states 11, 01, 10 and 00 according to distribution of a threshold voltage. First, a process of programming the LSB is the same as that of programming of the SLC. A memory cell having a state 11 is programmed to have a state A indicated by a dotted line according to the LSB.\nTo program the MSB, the memory controller transfers one-page data in the buffer memory 121 to the flash memory 130. Referring to FIG. 2, the memory cell in state A indicated by a dotted line is programmed to be in state 00 (program 1) or be in state 10 (program 2) according to the MSB. A memory cell in state 11 maintains state 11 or is programmed to be in state 01 (program 3) according to the MSB.\nReferring to FIG. 3, a memory cell is programmed to have any one of four states 11, 10, 00 and 01 according to distribution of a threshold voltage. First, a memory cell in state 11 maintains the state 11 or is programmed to be in state 10 (program 1) according to the LSB. Thereafter, the MSB is programmed. A memory cell in state 10 maintains the state 10 or is programmed to be in state 00 (program 2) according to the MSB. A memory cell in state 11 maintains the state 11 or is programmed to be in state 01 (program 3) according to the MSB.\nReferring again to FIG. 1, the memory system 100 uses the above-described method to program multi-bit data into the cell array 131 of the flash memory 130. That is, the LSB is programmed first in a memory cell, and then the MSB is programmed in the same memory cell.\nWhile the MSB is programmed after the programming of the LSB, a program fail may occur. In general, a possibility of a program fail is relatively low in the case of an SLC, while being relatively high in the case of an MLC because of a small threshold voltage margin of the cell. As the number of bits programmed in one MLC increases, the possibility of the program fail gradually increases.\nIn conventional MLC technology, a program fail occurs when the MSB is programmed after the LSB is normally programmed. In this case, a memory block including a fail cell is processed as a bad block. The fail cell cannot be restored even by an error correction circuit (ECC). Thus, in conventional MLC technology, because the memory block including a fail cell is processed as a bad block, even normal cells in the bad block become unavailable for an LSB programming operation."}
{"text": "The present disclosure generally relates to cartons or carriers for holding beverage containers or other types of articles. More specifically, the present disclosure relates to cartons that include reinforcement features for strengthening the carton."}
{"text": "The present invention relates to a method and a device for operating a battery, in particular a lithium ion battery, in a consumer (load).\nRechargeable batteries, in particular lithium ion batteries, are used in a multiplicity of products as energy stores. Currently available lithium ion batteries are challenging from a safety point of view. In the event of an accident, the battery can burn, emit toxic gases and dusts and even explode. The thermal runaway in the case of lithium ion batteries is primarily due to the accelerated oxidation of an organic electrolyte (frequently organic carbonates, such as, e.g., ethyl carbonate (EC), dimethyl carbonate (DMC), diethyl carbonate (DEC), etc.) at high temperatures. The decisive factor here is the exothermal self-acceleration, which a battery undergoes up to the eventual fire and the thermal runaway associated therewith. The fire is ultimately triggered by the release of oxygen by the metal oxide cathode material, which together with the organic of the electrolyte produces a combustible mixture. LiCoO2 and LiNiO2 are known for such instabilities.\nBy means of more stable cathode materials, such as, e.g., the so-called NCM or respectively NMC (Ni—Co—Mn based), the release of oxygen can be shifted to higher temperatures and the danger of a fire is thus reduced. Besides the composite oxide NCM, the so-called NCA (Ni—Co—Al based) is also used. Batteries which are NCA based on the cathode side are often used for high performance applications, such as, e.g., E-bikes. NCA based batteries are, however, more challenging from a safety perspective in comparison to other batteries that are based on oxidic cathode material. The primary reason for this is due to the high nickel content, which as a NCA represents a comparatively instable structure and can release oxygen at elevated temperatures. Said oxygen reacts strongly with the organic electrolyte and leads to a so-called thermal runaway, which in turn can lead to a fire or an explosion.\nIn order to prevent this, there is a need for a safe and reliable management of such a battery."}
{"text": "Cradles and auxiliary cradles for serving a secondary function such as supporting a receiver without disconnecting the circuit is generally shown in the patent to Clarke U.S. Pat. No. 1,391,370. With the advent of combining the receiver with the transmitter into a single handset, having pronounced protuberances with generally tapered cylindrical shapes, produced a large number of hangers and cradles, for example, see the U.S. Letters Pat. Nos. to B. H. King 2,339,413; S. G. Sears 2,443,329; J. J. Fleming 2,639,334; W. V. Fielder, Jr. 3,150,239 and D. T. Ross 3,538,267. Each of these devices were dependent upon the large protuberance of the receiver and transmitter extending laterally from the joining hand-piece and relied heavily upon the area of contact available. With the introduction of the Trimline.RTM. telephones designed by Bell Telephone Laboratories, of the type shown in the patent to S. W. Walden U.S. Pat. No. 3,557,322, the problems of controlling the proper usage of the telephone at the time of hang-up became more critical. These later telephones include a dial or tone buttons in the handset and damage to the handset resulted in a greater cost of repair and/or replacement. While the locking means discussed in the patent to Walden U.S. Pat. No. 3,557,322 is adequate when the device is utilized in the proper fashion, it functions well and its aesthetic appearance is pleasing to the eye, unfortunately, in many environments where there are teen-agers involved, where a patient in a hospital is attempting to hangup the telephone or in rough environments such as aboard ships or aircraft the limited area of contact at the receiver end has proven to be a deficiency."}
{"text": "With increasing in diameter and reducing in thickness of a size of a wafer for use in manufacture of a semiconductor device in recent years, a technology for sealing the wafer in a wafer level is demanded. Accordingly, in addition to a conventional transfer molding method using a solid type epoxy resin, a compression molding method using a liquid type epoxy resin has been suggested (Patent Literature 1).\nHowever, since the resin is flowed in a narrow portion in the transfer molding, there is concern of occurrence of wire transformation, and a short shot involved by an increase in sealing area is apt to occur. Further, in the compression molding method, it is difficult to finely control the molding range on an end face portion of a wafer, and it is not easy to optimize fluidity and properties when flowing a liquid sealing resin into a molding machine. Moreover, warpage of a molded wafer, which has been unquestioned, becomes an issue due to increase in diameter of a wafer size and reduction in film thickness of a wafer in recent years, and excellent wafer protection property is demanded. Accordingly, it has been demanded to develop a wafer mold material that can mold a wafer in a lump without causing problems such as a short shot on a wafer surface, and gives low-warping properties, as well as good wafer protection properties, good adhesion properties, high reliability, and good heat resistance after molding."}
{"text": "1. The Field of the Invention\nThis invention relates to fuel generation and, more particularly, to novel systems and methods for gasification of municipal solid waste and other feedstocks.\n2. The Background Art\nManufactured gas, or synthesis gas (syn-gas) as it is more often called today, comprises the unburned gasses (carbon monoxide and hydrogen gas) created by incomplete combustion of an organic raw material. Gasification (the process of generating manufactured gas) was discovered independently in both France and England. By 1850, gasification of coal had developed to where much of London was lit with manufactured gas, “town gas,” or “coal gas” as it was called. By 1920, many towns and cities throughout the United States supplied manufactured gas to their residents through local “gasworks.”\nFollowing 1930, as natural gas pipelines began to proliferate, low-cost natural gas displaced manufactured gas. The gasification industry was largely abandoned and forgotten. However, beginning with the oil embargo of the 1970's, there has been an almost continual increase in the cost of fuel. Accordingly, what is needed today is an efficient system, process, and gasifier that can receive a wide variety of inputs and efficiently produce a clean fuel."}
{"text": "Uncounted thousands of folding lounge chairs in a variety of shapes and sizes are sold each year for use on beaches, beside swimming pools and on patios and balconies. One of the most common and least expensive folding lounges consists of an adjustable tubular metal frame covered by a flexible plastic cover and sold in K-Mart and similar discount stores for less than $25.\nThe metal frame of these popular folding lounges is quite sturdy and always outlasts its plastic cover which in use soon tends to rip or tear or develop a permanent sag and thus fails to comfortably support a person lying or sitting on the lounge."}
{"text": "The benefits of regular exercise have been well established and accepted. However, due to time constraints, inclement weather, and other reasons, many people are prevented from activities such as participating in sports, walking, jogging, running, and swimming. As a result, a variety of exercise equipment has been developed. It is generally desirable to exercise a large number of different muscles so as to provide for balanced physical development, and to achieve optimum levels of exercise. It is further advantageous for exercise equipment to provide smooth and natural motion, thus avoiding significant jarring and strain that can damage both muscles and joints.\nWhile various exercise systems are known in the prior art, these systems suffer from a variety of shortcomings that limit their benefits and/or include unnecessary risks and undesirable features. For example, stationary bicycles are a popular exercise system in the prior art; however, these machines employ a sitting position and require the user's legs to move in a single, fixed, repetitive motion. Cross-country skiing exercise devices are also utilized to simulate the gliding motion of cross-country skiing. While cross-country skiing devices exercise more muscles than stationary bicycles, the substantially flat shuffling foot motion provided by the ski devices limits the range of motion of some of the muscles being exercised.\nTreadmills are still a further type of exercise device in the prior art. Treadmills allow natural walking or jogging motions; however, treadmills can enable significant impact loads to be transferred to the hips, knees, ankles, and other joints of the user, particularly when the treadmill is used to simulate running or jogging.\nAnother type of exercise device simulates stair climbing. Such devices can be composed of foot levers that are pivotally mounted to a frame at their forward ends and have foot-receiving pads at their rearward ends. The user pushes his/her feet down against the foot levers to simulate stair climbing. Resistance to the downward movement of the foot levers is provided by springs, fluid shock absorbers and/or other elements. These devices exercise more muscles than stationary bicycles; however, the rather limited range of up-and-down motion utilized does not exercise the user's leg muscles through a large range of motion. The substantially vertical reciprocating motion of such stair climbing exercise machines can result in substantial impact loads being applied to the hips, knees, ankles and other joints of the user. Further, the up and down reciprocating motion can induce a hyperextension of a user's knee.\nA relatively new class of exercise devices is capable of producing elliptical motion that better simulates the natural stride of a person. Elliptical motion is much more natural and analogous to running, jogging, and walking than the linear-type, back and forth motions produced by some prior art exercise equipment. However, these devices that create an elliptical motion are limited to analogizing to running, jogging, and walking motions.\nThe exercise devices of the prior art largely provide the user with a substantially fixed and limited range of motion that many users' find to be repetitious and uninteresting. It is well known that the more stimulating and enjoyable the experience of exercising is to a user, the longer and more frequently that user will exercise. The very advantage of the exercise equipment described above—the ability to use such equipment conveniently, in a relatively confined space, and in inclement weather—results in exercise devices that are relatively monotonous to use.\nAccordingly, a continuing need exists for an exercise device that provides an enjoyable stimulating experience and avoids being relatively monotonous to use. A need exists for an exercise device that enables a user to use his or her own creativity to define the exercise movements. There is also a need for an exercise device that enables the user to exercise muscles in a smooth natural manner, without applying undesirable impact loads to the user's joints. There is also a need for an exercise device that enables a user to improve his or her stability, coordination, rhythm and balance. It would be desirable for such an exercise device to be configured for convenient use in a relatively confined space even in inclement weather."}
{"text": "(1) Field of the Invention\nSubject invention is related to data transmission systems and more particularly to a scheme for correcting the crosstalk error between various groups or channels of a data transmission system.\n(2) Description of the Prior Art\nWhenever there is a transmission of data over multichannels of a transmission system, there is always some crosstalk between channels which contaminates the signals and thus poses all kinds of problems. Currently, crosstalk contamination of signals in a multichannel data transmission system is reduced by using only passive methods such as better cables or wider frequency separations. There is thus a need for an active method for correcting contaminated signals resulting from cross talk between multichannels of a transmission system."}
{"text": "Every year mobile phone users pay wireless providers billions of dollars for their services. By personalizing the mobile phone via contact lists, pictures, calendars, and individual settings, mobile providers have made the users become personally attached to these electronic devices.\nDue to these factors of convenience, familiarity and habit, it has been shown that users are more likely to use their mobile phones, even when given the choice of available free or low cost landline or voice over Internet Protocol (VoIP) communications. Indeed, nearly 3 out of every 10 households in the United States are reliant or almost completely dependent on mobile phones, and are willing to pay expensive mobile phone rates at times when the use of a mobile phone is unnecessary. As with using VoIP services, just the task of manually adding the name and phone number of each contact into the address book is both tedious and time consuming and the VoIP interface application is generally not as personal or familiar to a user as his own mobile phone.\nAnother trend is that in recent years, people have increasingly been spending more time using computers. Computers are readily accessible to users at various locations throughout the day, such as at an office, within one's home, and at recreational facilities, e.g. lounges, bars and coffee shops.\nIt would therefore be desirable to combine the accessibility of computers and the familiarity and personalization of mobile phones in order to establish a VoIP communications session.\nSkydeck, San Mateo, Calif., provides an application for accessing mobile phone data via a website. This application does not, however, generate images appearing similar to the mobile phone of the user, thereby making use of this application less personal and attractive than an actual mobile phone. Moreover, users of the Skydeck application need to interface with a fairly complex computer program, which can be confusing and challenging to many mobile phone users.\nDeviceAnywhere, San Mateo, Calif., provides a real-time platform by which mobile phones can be remotely tested. A replica of a mobile phone is displayed, and every input and output thereof is accessible via the Internet, so that buttons can be remotely pressed, liquid crystal displays (LCD) can be viewed, and tones can be heard. The service for performing this function is relatively expensive.                Skype Technologies S.A., Luxembourg, provides an application for supporting VoIP communication by means of a mobile phone.Users of such an application incur data transmission charges associated with using an Internet based VoIP application that interfaces with their mobile phone. They also incur charges for a local mobile phone call through their carrier, as well as a call connection fee and a subscription fee. Another disadvantage associated with this VoIP communication method is that a call cannot be directly placed by accessing the contact list of the mobile phone. In order to place a call, the user needs to manually enter the number of the intended recipient of the call.        \nIt is an object of the present invention to provide a method and system for establishing a computer assisted communications session.\nIt is an additional object of the present invention to provide a method and system for establishing a free or low cost communications session while interfacing with personalized settings and contact lists similar to those of the user's mobile phone.\nOther objects and advantages of the invention will become apparent as the description proceeds."}
{"text": "Conventionally used optical transmission systems transmit and receive optical signals among optical transmission devices having optical cross-connects (OXCs) by using the wavelength division multiplexing (WDM) method. In an optical transmission system of the WDM method, optical signals from each route are cross-connected in optical switches that make up the OXCs. Representative configurations of an optical transmission system of the WDM method include the broadcast and selector (BS) type in which input WDM signals are demultiplexed by a splitter and optical signals to be output are multiplexed by a wavelength selective switch (WSS), and the route and selector (RS) type in which input WDM signals are demultiplexed by a WSS and signals to be output are multiplexed by a WSS. In the BS type, in order to monitor input WDM signals by using an optical channel monitor (OCM), optical transmission devices have a 2×1 optical switch mounted at the stage prior to the OCM, and switch between monitored ports and monitor the wavelengths of input WDM signals. In contrast to this, in the RS type, in addition to monitoring the wavelengths of the input WDM signals, optical transmission devices also monitor optical signals that are multiplexed by a WSS and optical signals that are demultiplexed by a WSS. An RS type that uses WSSs for demultiplexing has less signal loss compared to a BS type that uses splitters for demultiplexing, and is advantageous for realizing a colorless, directionless, and contentionless (CDC) function.\nIn recent years, configurations in which the BS type and the RS type coexist within one optical transmission system have also been proposed. In an optical transmission system in which the BS type and the RS type coexist, optical transmission devices use a multi-input optical switch at the stage prior to an OCM in order to monitor optical signals from each port including also optical signals that are output. Thus, optical transmission devices are able to monitor the optical signals of each channel while individually selecting optical signals.\nJapanese Laid-open Patent Publication No. 2010-245993 is an example of related art."}
{"text": "The present invention relates to a liquid crystal display apparatus, and especially to an active matrix-type liquid crystal display apparatus.\nIn conventional active matrix-type liquid crystal display apparatus, a twisted nematic mode, a lateral electric field mode, MVA (Multidomain Vertical Alignment) mode, and so on which use nematic liquid crystal are used. Further, in those liquid crystal display apparatuses, there is a display method called a “hold type” display method, and in this display method, the same image is displayed during one display period of an image signal, that is, during one frame period.\nWhen dynamic images are displayed on the “hold type” crystal display apparatus, one image of the dynamic images which actually change moment by moment is held during one frame period. That is, although one point in the displayed image is displayed at the correct position at one moment, the point in the displayed image is different from the actual point in other moments. Thus, since a human perceives dynamic images by averaging the displayed images, the perceived images are not focused.\nA paper “Improving the Moving-Image Quality of TFT-LCDS” by K. Sueoka et al., IDRC '97, pp 203–206 (1998) discloses a technique in which, after the whole display panel has been scanned, a lighting device is turned on to eliminate the lack of focus due to the above averaging effect.\nHowever, in the above technique, since the lighting device is turned on after the whole liquid crystal panel has been scanned, and the response of the whole liquid crystal has been completed, the scanning period and the response time must be significantly shortened. Also, since the lighting period of the lighting device is short, the light strength must be increased in order to achieve the same brightness as that in a conventional liquid crystal display method. Thus, it is necessary to increase the light-tube current, which in turn decreases the lifetime of the lighting device."}
{"text": "The present invention comprises a new and distinct cultivar of Phalaenopsis Orchid, a bigeneric hybrid of Phalaenopsis×Phalaenopsis of the family Orchidaceae, and hereinafter referred to as the cultivar name, ‘SOGO F1582’. The genus Phalaenopsis is a member of the family Orchidaceae.\nPhalaenopsis comprise a group of bigeneric hybrids generally intermediate in character between the parent genera, which are suitable for cultivation in the home or greenhouse. The parent genera of Phalaenopsis are predominantly epiphytic or rock dwelling, and are native to tropical Asia, Malay Archipelago, and Oceania. The species typically have 2-ranked fleshy oblong or elliptic leaves affixed to a short central stem (monopodial growth), which vary in size from 5 to 8 inches to over 2 feet. The leaves may be entirely green or mottled with silver grey.\nPhalaenopsis orchid, often referred to ‘Moth Orchids’ in the horticultural trade, are frequently used to furnish cut flowers for the florist trade, or sold as flowering potted plants for home or interiorscape.\nPhalaenopsis produce upright racemes, often with many showy flowers, which open in succession beginning with the lowermost. The flowers possess three sepals, and three petals, the lateral ones being alike. The lowermost petal, called the labellum, is three lobed and is often more brightly colored than other flower segments. Flower colors are frequently various shades of pink, white and yellow.\nPhalaenopsis orchids are typically propagated from seeds. However, Phalaenopsis is capable of being asexually reproduced from offshoots, which frequently arise from the lower bracts of the inflorescence. The resulting plants are detached from the mother plant and may be planted in a suitable substrate.\n‘SOGO F1582’ is a product of a planned breeding program conducted by the inventor in Kaohsiung County, Taiwan, R.O.C. The objective of the breeding program is to create new uniform pot-type Phalaenopsis Orchid cultivars having attractive flower coloration. The inventor has addressed himself to Orchids breeding since 1985.\n‘SOGO F1582’ was discovered by the inventor from within the progeny of a cross-pollination of one Phalaenopsis Orchid, Phal. ‘Hilo Lip’ and one Phalaenopsis Orchid, Phal. ‘Schilleriana’ on February 1999, in a controlled environment in Kaohsiung County, Taiwan, R.O.C.\nAsexual propagation by tissue culture in a laboratory in Pingdong County, Taiwan, R.O.C. has been used to increase the number of plants for evaluation and has demonstrated in a controlled environment in Kaohsiung County, Taiwan, R.O.C. that the unique combination of characteristics as herein disclosed for the new Phalaenopsis Orchid are firmly fixed and are retained through successive generations of asexual reproduction."}
{"text": "The present invention relates to a coupling connecting an article handler to a source of pneumatic signals, and particularly to a rotary coupling for selectively blocking and restoring fluid communication between a rotating article handler and the stationary pneumatic signal source. More particularly, the present invention relates to a rotary coupling that couples a rotating article handler associated with a heat seal machine to the stationary source of pneumatic signals and that controls the communication of the pneumatic signals therebetween.\nFood processing equipment is frequently complex equipment having various translating or rotating moving parts. In addition, food processing equipment can employ pneumatic flow such as vacuum or forced air to either hold an article against an article handler or to separate the article from the article handler. What is needed is a rotary coupling that can couple a stationary source of pneumatic signals to a rotating portion of an article handler to allow fluid communication therebetween. Equipment manufacturers and operators of food processing equipment would appreciate food processing equipment having such a rotary coupling that is capable of communicating pneumatic flow through multiple channels and conduits, that is capable of restricting the communication of the pneumatic flow so that the pneumatic flow is available only during selected portions of the rotation cycle of the article handler, and that is capable of communicating both vacuum and forced air signals.\nAccording to the present invention, a rotary coupling is provided. The rotary coupling controls the communication of pneumatic signals from a stationary pneumatic signal source to an apparatus. The rotary coupling comprises a housing having an inlet in fluid communication with the pneumatic signal source and a rotating portion received in the housing for rotation with respect thereto. The rotating portion has an outlet in fluid communication with the apparatus. The inlet and the outlet are radially aligned relative to an axis of rotation of the rotating portion so that the inlet and the outlet are in fluid communication when the rotating portion is at an angular orientation at which the inlet and the outlet are coincident.\nIn preferred embodiments, the rotary coupling connects a rotating article handler associated, for example, with a heat seal machine to the source of pneumatic signals. The pneumatic signal source can illustratively include a vacuum source or a source of pressurized air. A rotating shaft, including a conduit in fluid communication with both the article handler and the rotary coupling, rotates the article handler. As the shaft rotates, the rotating article handler alternately lifts articles using suction force when the rotary coupling connects the article handler to the vacuum source and releases the articles when the rotary coupling removes the suction by disconnecting the article handler from the vacuum source and connecting the article handle to the source of pressurized air.\nThe rotary coupling in accordance with the present invention controls and \"programs\" the communication of pneumatic signals from the pneumatic signal source to the article handler. The rotary coupling establishes the sequencing of the pneumatic signals as well as the timing and duration of the signals. As the shaft rotates, the coupling selectively blocks and restores fluid communication between the article handler and the pneumatic signal source to communicate the pre-programmed pneumatic signals at each portion of the rotation cycle of the article handler. In addition, the coupling can be readily and easily adjusted to retard or advance the timing of the pneumatic signals. Also, the sequencing of signals can be changed entirely simply by changing a stationary pad inside of the coupling, completely \"reprogramming\" the sequence of pneumatic signals.\nThe rotary coupling includes a housing defining a cavity. The housing includes an end cap and the stationary pad which is supported by the end cap inside the cavity in a radially fixed orientation relative to the end cap for slight axial movement relative to the end cap. The stationary pad has an opening extending through the pad that is in fluid communication with an opening in the end cap. The opening in the end cap is in fluid communication with the source of pneumatic signals.\nThe rotary coupling also includes the shaft having a first end coupled to the rotating article handler of the equipment and a second end having an end portion rotatably received by the housing. The end portion includes a rotating pad supported in the cavity for slight axial movement relative to the shaft but in a radially fixed orientation relative to the shaft. The rotating pad thus rotates relative to the stationary pad.\nThe stationary pad has a first sealing face that is preferably very finely machined, ground flat, and then lapped to near perfect flatness. Likewise, the rotating pad has a second sealing face that is preferably very finely machined, ground flat, and then lapped to near perfect flatness. The shaft is received in the housing so that the first sealing face sealingly engages the second sealing face.\nSealing engagement of the first and second sealing faces is achieved in part by providing the first and second sealing faces with near perfect flatness as described above. In addition, the stationary and rotating pads are mounted so that the pads axially \"float\" between the end cap and the end portion of the main shaft. This float is achieved by providing the axially slidable connections between the stationary pad and the end cap and between the rotating pad and the end portion of the main shaft and by yieldably biasing the stationary pad away from the end cap and yieldably biasing the rotating pad away from the end portion of the main shaft so that the stationary and rotating pads are biased toward one another.\nThe stationary pad is yieldably biased away from the end cap by o-rings made from a resilient material positioned to lie between the stationary pad and the end cap. Likewise, the rotating pad is yieldably biased away from the end portion of the main shaft by o-rings made from a resilient material positioned to lie between the rotating pad and the end portion of the main shaft. These o-rings allow the stationary pad to move sightly in an axial direction relative to the end cap and the rotating pad to move slightly in the axial direction relative to the main shaft. The o-rings also allow the pads to adjust so that pads are not necessarily precisely parallel to the end cap or the end portion of the shaft, thus ensuring the sealing engagement between the sealing faces when opposing faces of the end portion and the end cap are not precisely parallel.\nThe rotating pad is formed to include an opening that extends through the rotating pad. The opening in the rotating pad is in fluid communication with the article handler through the conduit of the shaft.\nThe opening in the rotating pad and the opening in the stationary pad are radially aligned so that the openings coincide during a predetermined portion of each rotation cycle of the shaft and the rotating pad relative to the housing and the stationary pad. When the openings in the stationary and rotating pads are coincident, the openings are in fluid communication so that the article handler is in fluid communication with the pneumatic signal source. When the openings in the stationary and rotating pads are not coincident, the sealing engagement between the first and second sealing faces blocks the fluid communication between the openings in the rotating and stationary pads, blocking the fluid communication between the pneumatic signal source and the article handler. Thus, the positioning of the openings in the stationary and rotating pads operates to program the rotary coupling to communicate the pneumatic signal from the pneumatic signal source to the article handler when the article handler is at a predetermined angular orientation.\nThe rotating and stationary pads can additionally cooperate to define an arcuate channel in fluid communication with the openings in the rotating and stationary pads at selected portions of the rotation cycle. The channel operates to maintain the fluid communication between the openings in the rotating and stationary pads during portions of the rotation cycle at which the openings are not coincident. The channel, in effect, lengthens the portion of the rotation cycle of the shaft during which the openings in the rotating and stationary pads are in fluid communication, thereby lengthening the duration of the communication of the pneumatic signal from the pneumatic signal source to the article handler.\nLengthening the channel lengthens the portion of the rotation cycle during which the openings are in fluid communication. In addition, the channels can be formed so that the openings are not centered within the channel. Instead, the channel can be formed, for example, so that the fluid communication starts earlier in the rotation cycle and ends earlier in the rotation cycle than would occur if the opening was centered in the channel. Likewise, the channel can be formed so that the fluid communication is restored later and is blocked later in the cycle than would occur if the opening were centered in the channel. It can be seen, then, that the channel can be used to alter both the orientation of the shaft at which the openings in the rotating and stationary pads are in fluid communication and the extent of the rotation cycle during which the openings in the rotating and stationary pads are in fluid communication, thus varying the timing and duration of the communication of pneumatic signals between the pneumatic signal source and the article handler.\nIt is preferred that the channel is formed in the stationary pad, however channels can be formed in both the stationary and rotating pads or in only the rotating pad without altering the scope of the invention as presently perceived. It can be seen that if the channel is formed on the stationary pad, for example, the sequencing and duration of the communication of pneumatic signals between the article handler and the pneumatic signal source can be easily changed by replacing the stationary pad with a different stationary pad having a channel with a different configuration than the channel on the replaced stationary pad.\nAdditional objects, features, and advantages of the invention will become apparent to those skilled in the art upon consideration of the following detailed description of a preferred embodiment exemplifying the best mode of carrying out the invention as presently perceived."}
{"text": "Hangers for supporting articles from a vertical wall member and which provide vertical adjustment means are well known in the prior art. For example, U.S. Pat. No. 3,881,673, issued to Peterson on May 6, 1975 discloses a tensioning device designed primarily to support a filter bag in a baghouse gas filtering apparatus. The device comprises a pair of pivoted, scissor-like bars having ratchet teeth. The bars fit through a bore in a rigid support frame and the teeth engage the horizontal member of the support frame. Advancing the bars through the bore draws a suspensory means that is attached to a filter bag and to the advancing bars, thereby applying tension to the bag. U.S. Pat. No. 1,690,941 issued to Nickerson on Nov. 6, 1928 discloses a lamp socket securing means comprising a bracket member having a pair of spaced apart, parallel arms each having two co-planar portions that have a plurality of teeth extending from opposite sides thereof. The teeth interact with a spring clip for position adjustment. U.S. Pat. No. 952,053 issued to Tarleton on Mar. 15, 1910 discloses an adjustable pipe hanger comprising a stem having a plurality of notches or indentations along one edge that interacts with a spring loaded ball detent for vertical position adjustment. U.S. Pat. No. 178,733 issued to Carter et al. on June 13, 1876 discloses a chandelier height adjusting device which employs a pair of oppositely facing, pivotal cams having teeth-like projections secured to a vertically extending rod member that is spring loaded and to which the chandelier is attached. Another outer tube has a series of teeth-like projections on its inner surface that interact with the projections on the cam for position adjustment.\nThe present invention provides a novel method for position adjustment in that the means for adjustment is integral with the hanger. The prior art devices required either pivot means for disengaging the hanger member from the support or a manual manipulation of several interacting parts to adjust the supported article from one position to the next.\nAccordingly it is an object of the present invention to provide a simple, inexpensive hanger device for adjustably securing an article to a wall member.\nIt is a further object of the present invention to provide a hanger device in which one of either the bayonet member or the channel means is equipped with resiliently deformable projections that interact with projection receiving means on the other of either the bayonet member or the channel means to thereby allow position adjustment by simply advancing or retracting the bayonet member to the desired position.\nIt is a still further object of the present invention to provide a hanger device in which the force required to adjust the position of the article to be supported can be varied by relatively simple design changes to thereby permit use of the hanger device with a wide variety of articles of varying size and weight or for use in applications requiring child resistant adjustment.\nStill further objects of the invention will become apparent from the detailed description of the invention below."}
{"text": "The present invention relates to lamp holders in general and, more particularly, to a lamp holder that provides an accurate positioning of an arcuate lamp with compensation for thermal expansion of the lamp during operation.\nHigh intensity arcuate lamps are used in the semi-conductor industry as a light source for mask-to-wafer projection printers. The arcuate lamp must be accurately positioned with respect to the optical system of the projection printer. In addition, the lamp must be easily replaceable in the projection printer. Provisions must also be made for providing electrical power to the lamp and a flow of coolant fluid around the lamp during the lamp operation.\nIt is accordingly a general object of the present invention to provide an improved lamp holder for an arcuate lamp.\nIt is a specific object of the invention to provide a lamp holder for an arcuate lamp that accurately positions the lamp with respect to an optical system.\nIt is another object of the present invention to provide a lamp holder that is easily replaceable without sacrificing the positional accuracy of the lamp with respect to an optical system.\nIt is still another object of the present invention to provide a lamp holder for an arcuate lamp that compensates for thermally induced expansion of the lamp.\nIt is a feature of the invention that the lamp holder configuration provides a convenient handle to grip the holder without touching the arcuate lamp.\nIt is another feature of the invention that the lamp holder can be fabricated easily from readily available materials.\nIt is still another feature of the invention that replacement of coolant supply hoses, electrical power leads, and the lamp holder and lamp does not affect the positional location of the lamp with respect to an optical system."}
{"text": "Means for separating substances from fluids have been employed in a wide variety of fluidic applications including brewing of alcoholic beverages, hydraulic systems, fuel systems, and engine lubrication systems.\nIn some applications, fluid filters may include fuel/water separators or sediment pots. In a fuel/water separator, water and sediments such as dirt, sand and grit are separated from a fuel/water mixture, in order to prevent damage to downstream engine components. Fluid filters which remove water will tend to accumulate the separated water and sediment by gravity at the bottom of the housing. The water should eventually be removed from the housing. Some models of liquid filters incorporate a mechanism to remove the water from the housing by using a pipe plug or a petcock. Often, the device begins to drain fluid as soon as it becomes unseated or loosened. The leaking fluid may run onto the operator's hand and down the arm as the device is turned the additional revolutions to the open, or separated, position. When the mechanism is to be closed, the operator is exposed to the fluid as the threaded shaft or plug is rotated several turns before it becomes closed or seated.\nIn other applications, for example, brewing of beer, wine, and other bottle-conditioned alcoholic beverages, i.e., beverages which are fermented, aged and naturally carbonated in the bottle may require the removal or addition of substances from the product. In the fermentation process used to prepare such beverages, yeast causes sugars in the liquid to ferment into carbon dioxide (CO2) and ethyl alcohol (C2H6O). During fermentation, the carbon dioxide may cause a relatively high pressure to build-up, in some cases up to 12 atmospheres. There are also various unwanted by-products of fermentation that vary according to the chemical composition of the liquid and the rate and manner of fermentation. These by-products either dissolve in the wine or precipitate as sediments. Various methods have been employed to remove excess carbon dioxide and unwanted by-products from the beverage or to add products to the beverage.\nIn the case of champagne making, for example, the bottle may be inverted or turned upside down to allow the sediment to collect in the tip region of the neck of the bottle. The bottle neck may then be placed in a freezing brine solution until the liquid in the tip region is frozen solid. The bottle may then be warmed slightly to loosen the frozen sediment plug, after which the bottle cap is removed and the pressure of the natural carbonation blows the sediment plug out of the bottle. Then the bottle may be recapped. This method, however, is complicated and time-consuming.\nFor the foregoing reasons, it is desirable to have a superior apparatus, method and system for venting and removing sediment from a liquid than what is currently available."}
{"text": "Hyperspectral imaging is a form of spectral imaging wherein information from across the electromagnetic spectrum is collected in many narrow spectral bands and processed. From the different spectral images that are collected, information of the objects that are imaged can be derived. For example, as certain objects leave unique spectral signatures in images which may even depend on the status of the object, information obtained by multi-spectral imaging can provide information regarding the presence and/or status of objects in a region that is imaged. After selection of a spectral range that will be imaged, as spectral images in this complete spectral range can be acquired, one does not need to have detailed prior knowledge of the objects, and post-processing may allow to obtain all available information. Whereas originally hyperspectral remote sensing was mainly used for mining and geology, other applications such as ecology, agriculture and surveillance also make use of the imaging technique.\nIt is known to use photogrammetric techniques to infer three-dimensional information, in particular elevation information, from the acquired two-dimensional images. An example of such a technique is disclosed in Alsadik, B. S., Gerke, M., & Vosselman, G. (2012), “Optimal Camera Network Design For 3D Modeling Of Cultural Heritage”, ISPRS Annals of the Photogrammetry, Remote Sensing and Spatial Information Sciences, I-3, 7-12.\nSome agricultural and ecological applications are known wherein hyperspectral remote sensing is used, e.g. for monitoring the development and health of crops, grape variety detection, monitoring individual forest canopies, detection of the chemical composition of plants as well as early detection of disease outbreaks, monitoring of impact of pollution and other environmental factors, etc. are some of the agricultural applications of interest. Hyperspectral imaging also is used for studies of inland and coastal waters for detecting biophysical properties. In mineralogy, detection of valuable minerals such as gold or diamonds can be performed using hyperspectral sensing, but also detection of oil and gas leakage from pipelines and natural wells are envisaged. Detection of soil composition on earth or even at other planets, asteroids or comets also are possible applications of hyperspectral imaging. In surveillance, hyperspectral imaging can for example be performed for detection of living creatures.\nInternational patent application publication WO 2011/073430 A1, in the name of the present applicant, discloses a sensing device for obtaining geometric referenced multi-spectral image data of a region of interest in relative movement with respect to the sensing device. The sensing device comprises a first two dimensional sensor element and a spectral filter. The spectral filter and the first sensor element are arranged for obtaining spectral information at a first wavelength or wavelength range using a part of the first sensor element and for obtaining spectral information at a second wavelength or wavelength range using another part of the first sensor element. As a result of this arrangement, different parts of a single image acquired with the first sensor will represent the imaged scenery as seen is radiation of a different respective wavelength.\nTo date, there is no satisfactory way to apply the aforementioned photogrammetric techniques directly to multi-spectral or hyperspectral images such as those acquired by means of the first sensor of WO 2011/073430 A1. Nevertheless, full photogrammetric information about the acquired multi-spectral or hyperspectral images is essential to correctly stitch together an image of the scanned scenery in any particular wavelength band.\nIn WO 2011/073430 A1, this problem is solved by providing a second sensor on the same substrate, which simultaneously acquires panchromatic images of the area imaged by the first sensor. The panchromatic images are used to perform 3D modeling, and the 3D model information is subsequently transposed to the composite images of the scenery in the different wavelength bands that are covered by the multi-spectral or hyperspectral first sensor.\nUnpublished international patent application no. PCT/EP2015/065523, in the name of the present applicant, describes a sensing device for obtaining geometric referenced multi-spectral image data of a region of interest in relative movement with respect to the sensing device, the sensing device comprising: at least a first two-dimensional sensor element, the sensing device being adapted for obtaining subsequent multi-spectral images during said relative motion of the region of interest with respect to the sensing device thus providing distinct spectral information for different parts of a region of interest using the first sensor element; a second two-dimensional sensor element, the sensing device being adapted for providing, using the second sensor element, an image of the region of interest for generating geometric referencing information to be coupled to the distinct spectral information; the first two-dimensional sensor element being operable to capture a first sequence of frames at a first frame rate and the second two-dimensional sensor element being operable to capture a second sequence of frames at a second frame rate; wherein the first frame rate is higher than the second frame rate; and wherein the sensing device further comprises a processor configured to generate intermediate geometric referencing information to be coupled to frames of said first sequence of frames for which no synchronous frame from said second sequence of frames is available, said intermediate geometric referencing information being derived from one or more temporally similar frames from said second sequence of frames. In this way, the system of PCT/EP2015/065523 is able to reduce the number of panchromatic images required to perform geometric referencing, relative to the system of WO 2011/073430 A1, by means of interpolation or extrapolation. However, the system of PCT/EP2015/065523 does not change the assumption that panchromatic images remain necessary.\nThe article by inventor of the present invention A. Sima et al., “Semi-Automated Registration of Close-Range Hyperspectral Scans Using Oriented Digital Camera Imagery And A 3D Model”, The Photogrammetric Record, vo. 29 (2104), no. 145, discloses a semi-automated method for registering terrestrial panoramic hyperspectral imagery with LIDAR models and conventional digital photography. The method relies on finding corresponding points between images acquired in significantly different parts of the electromagnetic spectrum, from different viewpoints, and with different spatial resolution and geometric projections. The article recognizes that the number of homologous points that can be matched between a hyperspectral band and a covering (panchromatic) digital camera image varies with the wavelength, and proposes a technique that relies on a judicial selection of spectral bands in which a maximum number of points can be matched.\nUnpublished international patent application no. PCT/EP2015/065524, in the name of the present applicant, describes a method for performing photogrammetric 3D reconstruction of objects imaged in a sequence of images, the images containing distinct regions representing imaged radiation in respective distinct wavelengths, the method comprising: selecting a plurality of subsets from the sequence of images, each one of the plurality of subset containing a plurality of images, each image of which represents a field of view that overlaps with a field of view of at least one other image in the same subset; generating a set of intermediate 3D models by performing photogrammetric 3D reconstruction on the images in respective ones of the subsets; and recombining the intermediate 3D models from the set of 3D models into a combined 3D model. However, this method still requires a significant amount of processing.\nInternational patent application publication no. WO 2014/151746 A2 in the name of URC Ventures Inc. discloses techniques for analyzing images acquired via mobile devices in various ways, including to estimate measurements for one or more attributes of one or more objects in the images. For example, the described techniques are stated to be usable to measure the volume of a stockpile of material or other large object, based on images acquired via a mobile device that is carried by a human user as he or she passes around some or all of the object. During the acquisition of a series of digital images of an object of interest, various types of user feedback may be provided to a human user operator of the mobile device, and particular images may be selected for further analysis in various manners. Furthermore, the calculation of object volume and/or other determined object information may include generating and manipulating a computer model or other representation of the object from selected images. The techniques disclosed in this article are not concerned with images that represent radiation in distinct wavelengths.\nThe article by Lichun Li et al., “A new navigation approach of terrain contour matching based on 3-D terrain reconstruction from onboard image sequence”, Science China Technological Sciences, vol. 35 (2010), no. 5, 1176-1183, discloses a passive navigation method of terrain contour matching by reconstructing the 3-D terrain from the image sequence acquired by an onboard camera. To achieve automation and simultaneity of the image sequence processing for navigation, a correspondence registration method based on control points tracking is proposed which tracks the sparse control points through the whole image sequence and uses them as correspondence in the relation geometry solution. A key frame selection method based on the images overlapping ratio and intersecting angles is explored, and then the requirement for the camera system configuration is provided. The proposed method also includes an optimal local homography estimating algorithm according to the control points, which helps correctly predict points to be matched and their corresponding speed. Consequently, the real-time 3-D terrain of the trajectory thus reconstructed is matched with the referenced terrain map, and the result provides navigating information. The techniques disclosed in this article are not concerned with images that represent radiation in distinct wavelengths.\nThe article by Turner Darren et al., “Direct Georeferencing of Ultrahigh-Resolution UAV Imagery”, IEEE Transactions on Geoscience and Remote Sensing, vol. 52 (2014), no. 5, 2738-2745, discloses a concept for a camera-global positioning system (GPS) module that allows the synchronization of camera exposure with an airframe's position as recorded by a GPS with 10-20-cm accuracy. Lever arm corrections were applied to the camera positions to account for the positional difference between the GPS antenna and the camera center. Image selection algorithms were implemented to eliminate blurry images and images with excessive overlap. This study compared three different software methods (Photoscan, Pix4D web service, and an in-house Bundler method), evaluating each based on processing time, ease of use, and the spatial accuracy of the final mosaic produced. The article mentions the need to remove some of these overlapping images to improve processing efficiency, which is achieved by reading the image location information from image EXIF headers. An image overlap of around 80% is stated to yield the best results for SfM-based image processing algorithms. Image subsetting is achieved by computing the distance between consecutive images. The techniques disclosed in this article are not concerned with images that represent radiation in distinct wavelengths.\nInternational patent application publication no. WO 2014/031284 A1 in the name of Visual Intelligence LP, discloses an imaging sensor system comprising: a mount unit affixed to a vehicle or a platform and having at least three imaging sensors disposed within the mount unit, wherein a first, second and third imaging sensor each has a focal axis passing through an aperture in the mount unit, wherein the first image sensor generates a first image area of the target area comprising a first array of pixels, wherein the second image sensor generates a second image area of the target area comprising a first array of pixels, wherein the first and second imaging sensors are offset to have a first image overlap area in the target area, wherein the first sensors image data bisects the second sensors image data in the first image overlap area. Various embodiments are described. The camera array assembly is configured such that adjoining borders of the relevant image areas overlap slightly. While this document suggests ortho-rectifying the sequence of overlapping images using standard photogrammetry techniques to produce an orthomap in which each pixel has an unique latitude and longitude coordinate and a unique elevation coordinate, it does not address dealing with images acquired in different spectral ranges.\nInternational patent application publication no. WO 2004/027348 A2, in the name of M7 Visual Intelligence LP, discloses a method to calibrate an on-board remote sensing system using a self-locking travel pattern and target remote sensing data. The self-locking travel pattern includes a number of parallel travel lines having overlapping swath widths between adjacent travel lines. The overlapping swath widths are used to determine the boresight angles and range offset of the remote sensor device. In addition, the method can be used to generate estimated horizontal and vertical displacement errors. These estimated errors can be used as correction factors for the range offset and boresight angles.\nThere is a need for an alternative way of geometrically referencing multi-spectral or hyperspectral images that does not rely on simultaneously acquired panchromatic images and that has limited computational complexity."}
{"text": "This invention relates generally to spectrophotometry, and more particularly concerns advancements in optical coupling to monochromators and between monochromator sections.\nTo facilitate understanding the invention and distinguishing if from prior art, it is necessary to define certain terms. A monochromator can be viewed as an optical filter. A double monochromator, then, is a system which causes the radiation to be twice filtered. The advantages of the double monochromator over the single monochromator are improved spectral purity and resolution. In this context, a monochromator stage is defined as a set of optical elements, including dispersing means such as a grating, necessary to receive light from one slit, render it incident to the dispersing means, and finally pass it through another slit. A double monochromator then can be either a system which passes the radiation once through two separate monochromator stages or a system which passes the radiation twice through a single monochromator state. The latter type are frequently referred to as double-passing monochromators. The advantage of the double-passing monochromator over the two stage monochromator is largely economic. It is generally a more compact system and results in fewer essential expensive optical and mechanical components. Also, the engineering problems of tracking the two monochromator sections to the desired accuracy are substantially reduced.\nA further complication exists in that in either type of double monochromator, the optics can be arranged so that the dispersion of the second section either adds to or tends to subtract from the dispersion of the first section. In the subtractive dispersion arrangement, because of the bandwidth limit imposed by the intermediate slit, the resultant total dispersion is equal to that of a single monochromator whereas the total resultant dispersion of the additive arrangement is twice that of the single monochromator. Since the amount of light flux passed for a given resolution is proportional to dispersion, for most applications the preferred embodiment is additive dispersion.\nIn the past, the achievement of desired monochromator efficiency has been limited by a number of factors. For example, in the case of additive dispersion double-pass monochromators, with intermediate slits located at the monochromator side of the plane or planes of the entrance and exit slits, bands of unwanted radiation are efficiently transmitted, requiring filters or selective detectors for avoiding or reducing response to such radiation bands and consequent spurious and misleading signals. See for example U.S. Pat. No. 2,922,331 to Fastie.\nExternal slit-coupling arrangements (and some internal ones as well) have produced serious mismatch of the curvature of the exit slit of the first section as imaged upon the entrance slit of the second section. Consequently, all such arrangements described have used very short slits so that resolution is not seriously deteriorated, with the result that the \"light grasp\" (or \"throughput\", or \"etendue\"), i.e. the amount of light flux that can be passed for a given resolution, is seriously limited. For illustrations of these arrangements, see U.S. Pat. Nos. 2,922,331 and 3,567,323."}
{"text": "FIGS. 1A and 1B schematically show light 5 entering from the left into a diffractive optical element 10 comprising a diffractive layer 11 having a repeated diffraction spacing, d. The light 5 is transmitted through the interior of the diffractive optical element 10 by total internal reflection and is diffracted by the diffractive layer 11. When diffracted, some of the light 5 is out-coupled from the diffractive optical element 10, as represented by arrows 15. The angles θn of the out-coupled light for each diffracted order n are determined by the wavelength λ of the light source and the repeated diffraction spacing d of the diffractive layer 11 according to the well-known equation:d(sin θn+sin θi)=nλ  [Eqn. 1]where θi is the angle of incident light, n is an integer, θn is the angle of the diffracted light and λ is the wavelength of the light. As shown in FIG. 1A, incoming light 5 entering the diffractive optical element 10 at smaller angles generally travels a lesser distance by each total internal reflection before being out-coupled by diffraction, so that the out-coupled light 15 is mostly concentrated near to where the light has entered the diffractive optical element 10, fading out rapidly from left to right. On the other hand, as shown in FIG. 1B, incoming light 5 entering the diffractive optical element 10 at larger angles can travel a greater distance by each total internal reflection before being out-coupled by diffraction, so that the out-coupled light 15 is mostly concentrated far from where the light has entered the diffractive optical element 10. Similar considerations would apply in a left-right mirror image if the incoming light 5 were instead to enter the diffractive optical element 10 from the right in FIGS. 1A and 1B, rather than from the left.\nFIG. 1C schematically shows a pair of such out-coupling diffractive optical elements 10L, 10R, each of which comprises a respective diffractive layer 11L, 11R. The out-coupling diffractive optical elements 10L, 10R are arranged in optical communication with an in-coupling diffractive optical element 30 comprising a similar diffractive layer 31 to that of the out-coupling diffractive optical elements 10L, 10R and having the same repeated diffraction spacing d. The in-coupling diffractive optical element 30 receives incoming light 5, which is diffracted by diffractive layer 31 and transmitted by total internal reflection to the out-coupling diffractive optical elements 10L, 10R.\nFIG. 1C also schematically shows eyeballs 81, 82 of a viewer gazing at the out-coupling diffractive optical elements 10L, 10R. As may be understood from FIG. 1A, for a viewer looking to the right, as shown in FIG. 1C, the out-coupled light 15 will therefore have lowered brightness in the region A, whereas light in the region B will fall outside the gaze of the viewer and be wasted. On the other hand, as may also be understood from FIG. 1B, for a viewer looking to the left, the out-coupled light 15 will similarly have lowered brightness in the region B and will fall outside the gaze of the viewer in region A and be wasted.\nFIG. 2A schematically shows a pair of components 1, 2 of an apparatus for out-coupling polychromatic light, for use, for example, in a binocular near-eye display. By polychromatic light is meant light of at least two different wavelengths. The components 1, 2 each have left and right halves which are mirror images of each other, configured to out-couple light to a pair of eyes. Component 1 therefore comprises a pair of out-coupling diffractive optical elements 10L, 10R, a pair of corresponding in-coupling diffractive optical elements 30L, 30R, and a pair of intermediate optical elements 51, 52, which respectively direct light from the in-coupling diffractive optical element 30L to the out-coupling diffractive optical element 10L and from the in-coupling diffractive optical element 30R to the out-coupling diffractive optical element 10R. Component 2 similarly comprises a pair of out-coupling diffractive optical elements 20L, 20R, a pair of corresponding in-coupling diffractive optical elements 40L, 40R, and a pair of intermediate optical elements 53, 54, which respectively direct light from the in-coupling diffractive optical element 40L to the out-coupling diffractive optical element 20L and from the in-coupling diffractive optical element 40R to the out-coupling diffractive optical element 20R.\nComponents 1 and 2 differ from each other in that the in- and out-coupling diffractive optical elements 30L, 30R, 10L, 10R of component 1 have a first repeated diffraction spacing, d1, whereas the in- and out-coupling diffractive optical elements 40L, 40R, 20L, 20R of component 2 have a second repeated diffraction spacing, d2, which is different from the first spacing, d1. Components 1 and 2 can therefore provide respective channels of optimized efficiency for diffracting light in two different wavelength bands. For example, component 1 may provide a channel for red light and component 2 may provide a channel for green-blue light. Thus if components 1 and 2 are superposed one on top of the other, as is schematically represented in FIG. 2B, and if the optical elements of each component are carefully aligned, polychromatic light from a single display can be projected into the in-coupling diffractive optical elements 30L, 30R, 40L, 40R. Light in two different wavelength bands will then be out-coupled from elements 10L and 20L and will re-combine to provide polychromatic light to a viewer's left eye, whilst light in two different wavelength bands will also be out-coupled from elements 10R and 20R and re-combine to provide polychromatic light to the viewer's right eye. A similar arrangement can be used in a monocular near-eye display to out-couple polychromatic light to a single eye by using only a left or right half of each of the components 1 and 2.\nThe components 1, 2 may each be understood as being similar in construction and function to the in- and out-coupling diffractive optical elements described above in relation to FIG. 1C. Therefore, they suffer from the same problems as were explained above in relation to FIG. 1C. Typically, for example, if the components 1, 2 are incorporated into a binocular near-eye display, they will provide a field of view of less than about 40 degrees. The “eye box” can be increased by scaling the size of all of the optical elements, starting with the display, but this is undesirable, from the point of view not only of cost, but also of wearability. On the other hand, a much wider field of view would be highly desirable, considering that the natural field of view of a healthy human can extend beyond 180 degrees in the horizontal direction."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an apparatus for injecting drugs into a patient and, more specifically, to an automatic injection device that utilizes a flexible membrane assembly and a propellant to automatically inject a single dose of a drug into a patient when the propellant is ignited.\n2. Background of the Invention\nPresently available automatic injection systems are typically large and bulky devices requiring skill and training in order to properly operate the device. Accordingly, such devices are not amenable to applications for which rapid administration of either a single drug or a combination of different compounds in rapid succession is desirable. Such systems are commonly used for mass vaccinations and self-administration of medications--such as for injection of insulin in diabetic patients, epinephrine in anaphylactic shock, and antimigraine drugs (e.g., sumatriptan). Devices for self-administration typically require cleaning after each use. Automated injectors used for mass vaccinations, as in military applications, typically use the same injection head for all patients and often require a supply of compressed gas.\nAutomatic injection devices are also presently in use for parenteral administration of medications under field conditions, where injections with a syringe are either impractical, require undue attention or training, or are painful. Examples of such situation s are administration of chemical warfare antidotes (e.g., atropine) and treatment of anaphylactic shock (e.g., epinephrine). In both instances, drug administration must be rapid, simple and reliable and must minimize the possibility of human error on selection and administration of drug doses. One such automatic injection device used by the military comprises an elongate tube with a needle protruding from one end. The needle is rapidly and forcefully inserted into a muscle, resulting in the painful injection of a drug contained within the tube. Upon removal of the needle, however, the needle remains exposed and thus requires proper handling to prevent contact by and/or injury to other persons.\nFor mass vaccination programs, reduction of the delivery costs associated with immunizations presents a big opportunity to boost immunization coverage rates. One way to decrease costs is to reduce the number of trained personnel required to complete an immunization series. Frequently, patient compliance is a major factor limiting effectiveness of mass vaccination programs--children, parents, and health care providers are less willing to accept the increased numbers of separate injections. Thus, it is desirable to deliver as many vaccines as possible during a single, quick patient visit. Mixing of vaccines in one syringe is often contraindicated due to potential incompatibilities and a reduced potency of the mixed vaccines. However, vaccines are typically safe and effective when administered at the same time at different application sites, such as with a single, multiple-site injection. Therefore, there is a clear need to simultaneously inject several vaccines into a patient with each vaccine being delivered at a discrete injection site. Due to problems associated with typical needle injections, including transmission of disease between patients, hazards to health-care workers, and expensive handling and disposal procedures, a method and apparatus for injecting a patient without using needles or by utilizing retractable needles is desirable.\nU.S. Pat. No. 5,167,625 to Jacobsen et al. discloses an implantable drug delivery system which includes a housing having a plurality of compartments, each of which includes a flexible drug containment sack. A cover is disposed over the compartment openings to prevent the inadvertent release of drug formulations contained in the sacks. Gas generating elements are disposed in the compartments and are responsive to an initiation signal for producing gas to force the sacks and drug formulations out of the housing. Such a device, however, is designed for internal use, requiring surgery for implantation.\nA portable, hand-held injection device is disclosed in U.S. Pat. No. 5,616,132 to Newman, and includes a diaphragm mounted within a housing having a needle associated therewith that moves in response to movement of the diaphragm when pressurized gas is released into the housing. Each device, however, is a completely self-contained unit that would be relatively expensive to manufacture, and difficult, if not impossible, to reuse. Accordingly, such a device would be impractical for mass injection scenarios as previously discussed.\nThus, it would be advantageous to provide prefilled, interchangeable and disposable injection devices that are inexpensive, easy to use, and safe for both the patient and medical personnel. It would also be advantageous to provide a prefilled injection device that reduces the risk associated with cross-infection, improves compliance--especially in children--improves the effectiveness of vaccination programs, does not require special sharps handling for disposal, and is a self-contained unit that presents a single item with few contaminated components for disposal."}
{"text": "1. Field of the Invention\nThe present invention relates to a magnetic powder transporting device, and more particularly to such device capable of preventing the magnetic powder scattering.\n2. Description of the Prior Art\nAs an example of the magnetic powder transporting device there is already known a sleeve-type developing unit for powdered toner for use in the electrophotographic copier, for example disclosed in the U.S. Pat. No. 4,091,765. The abovementioned device is composed of a magnet roller provided on the peripheral surface thereof with magnets and a non-magnetic sleeve covering said peripheral surface, wherein said magnet roller and sleeve being maintained in mutual rotation to transport the powdered toner to a developing area where an electrostatic latent image on a photosensitive member is rendered visible. In said device the toner transportation is achieved either indirectly by iron powder carrier in case of two-component toner consisting of iron powder and toner particles, or directly by the toner particles in case of one-component toner consisting solely of toner particles having iron cores. Said toner transportation is conducted to replenish new toner to the developing area in which the toner is consumed.\nThe above-mentioned sleeve-type developing unit is advantageous in the significantly smaller weight and dimension in comparison with the cascade developing unit for example disclosed in the U.S. Pat. No. 2,573,881, but is associated with a drawback of causing toner scattering from the lateral ends of the developing area, thus severely soiling the interior of the copier. In the above-mentioned cascade developing unit, the toner particles transported to the developing area are subjected merely to gravity and not influenced by any external force towards the lateral ends of the developing area. In contrast, in the sleeve-type developing unit, the toner particles supported on the sleeve in spicate forms are spread, upon contact with the surface of the photosensitive member, over said surface, thus reaching the lateral ends of the developing area. Said toner spreading, functioning as a partial force toward the lateral ends of the developing area, naturally affects the toner particles present in said lateral end regions where the magnetic field ends, thereby releasing said toner particles from the magnetic field and causing the toner scattering.\nThe toner scattering is also caused by a fact that, although the magnetic field around the periphery of the sleeve is formed between the adjacent magnetic poles and is generally parallel to the plane perpendicular to the axis of the sleeve whereby the toner particles are scarcely subjected to axial force, the magnetic field at the lateral ends thereof is directed along the axis of the sleeve, thus causing the displacement of toner particles in said direction and facilitating the toner scattering.\nIn the conventional device as shown in FIGS. 1 and 2 respectively giving a partial view of the sleeve-type developing unit and a partial perspective view of the magnet roller, the toner particles T captured and transported on the surface of a sleeve 2 by a magnet roller 1 are shaped in almost vertical spicate forms over the approximately entire surface of said sleeve 2 until they come into contact with the surface of the photosensitive member 3, but said spicate forms are spread in lateral directions upon such contact. The magnetic field formed on the magnet roller 1 is mostly parallel to the plane perpendicular to the axis of the magnet roller as shown by the arrow A in FIGS. 1 and 2, but at the end portions of the roller there are formed magnetic fields in the axial direction as shown by the line B. (Although the magnetic field at the end portions is in fact complicated, there are only shown the lines in the axial direction for simplicity.) For this reason the toner particles T spread in the above-explained manner are subjected, at the end portions of the magnet roller 1, to the axial force along the magnetic field B, thus easily leaving the end faces of the roller and causing the toner scattering."}
{"text": "1. Field of the Invention\nThis invention relates to a pharmaceutical suspension. More particularly, the instant invention is concerned with a suspension of titanium dioxide for use in admixing with gelatin solutions from which opaque empty gelatin capsules can be manufactured which are suitable as containers for unit doses of pharmaceutically active agents.\n2. Prior Art\nAlmost all of the billions of empty gelatin capsules used each year in the world are manufactured by dipping highly polished, but lubricated, forming pins into a warm viscous gelatin solution. The forming pin is slowly withdrawn from the solution with a thin film of gelatin adhering thereto. The gelatin film is dried on the forming pin until the moisture content of the film is somewhere in the neighborhood of 13-16 percent, preferably about 14-15 percent. A special stripping tool is used to remove the gelatin film from the forming pin and the edges of the open end of the gelatin are trimmed to provide one of the two pieces of an empty gelatin capsule. Separate forming pins of different diameters are used to form each of the two pieces of an empty capsule so that, when properly trimmed, the two pieces will telescope together to make a package which can be filled with various physiologically active agents to provide a pharmaceutical unit dosage form.\nGelatin is a material which has been used since the late nineteenth century for making capsules for use in holding unit dosage forms of various active agents.\nA simple solution (actually a colloidal solution) of gelatin can be used to make capsules in the manner described briefly above. It is preferred, however, to include preservatives in the gelatin solution because in such solution gelatin is a highly desirable substrate for microbial growth. And when the organisms are gas formers, little bubbles appear in the solution which are not easy to remove.\nEarly capsules were made without coloring and, if the gelatin was of good quality, such capsules would be essentially transparent. Later on various approved soluble dyes were added to the gelatin solution to make capsules of various colors. Even so, such capsules, colored or plain, came to be known as clear capsules.\nAs the demand for empty gelatin capsules increased, a need arose for a wider range of colors and hues for use in identifying different medicaments. A group of opaque capsules were developed to meet this need, and titanium dioxide was adopted as the opaquing agent because of its lack of toxicity and high opaquing properties.\nOpaquing grade titanium dioxide must be finely divided particles, preferably in the micron range. Generally, particles this small are difficult to wet, and titanium dioxide is no exception. It was soon found that it was impractical to mix finely-divided titanium dioxide powder directly into a viscous gelatin solution. Consequently, paste formulations of titanium dioxide were developed which could be mixed into the viscous gelatin solution with only moderate difficulty. However, these titanium pastes had to be prepared very nearly simultaneously with the preparation of the gelatin solution as the pastes tended to coalesce and harden within a short time.\nThen it was discovered that the addition of glycerin to the paste formulation would substantially retard the hardening and make it possible to prepare such a paste many days ahead of its use.\nSometime later a suspension of titanium dioxide was developed which was suitable for admixing with a solution of gelatin for the manufacture of opaque gelatin capsules. And, for many years this suspension proved satisfactory for its purpose. It consisted essentially of titanium dioxide powder, glycerin, water, and the preservatives methylparaben, propylparaben. Later chloroform and alcohol S.D.A. No. 20 were added.\nThen a few months ago it became impossible to obtain adequate amounts of the grade of titanium dioxide which had been used for years in preparing the suspension. A new supplier of titanium dioxide was found and a different grade of titanium dioxide was substituted for that which had been used in the past. With the new grade of titanium dioxide it was no longer possible to prepare a suitable suspension having a viscosity within the range required for effective additions to gelatin solutions. And, when a mixture of the titanium dioxide and the gelatin solution was accomplished, the capsules produced therefrom were grainy and reflected light unevenly. It was learned that, even though an analysis of the particle size of the two grades indicated a similar range of sizes, an average of about 1 micron, one grade, the longused quality, was specified as an oil dispersible material, whereas the newer material was categorized as a water dispersible grade. This classification was supported by the fact that a 25 percent aqueous suspension of the so-called oil dispersible grade had a viscosity of between 1,400 and 2,200 centipoises at 25.degree. C., and a similarly prepared aqueous suspension of the water dispersible grade had a viscosity of about 3 centipoises at 25.degree. C. It was also found that the water dispersible grade of titanium dioxide had a zeta potential of about -35 millivolts whereas the oil dispersible material had exhibited a zeta potential of -17 to -27 millivolts.\nAccordingly, it is an object of this invention to provide a pharmaceutical suspension of titanium dioxide of either oil or water dispersible grades that can be admixed with gelatin solutions for manufacturing empty opaque capsules of consistent quality.\nAnother object of this invention is to provide a process for preparing a pharmaceutical suspension of titanium dioxide utilizing either oil or water dispersible titanium dioxide, said suspension being for use in admixing with gelatin solutions for manufacturing empty opaque capsules of consistent quality.\nStill another object of this invention is to provide a pharmaceutical suspension of titanium dioxide of whatever source or grade that resists caking and settling of the dispersed phase and which, after dead storage for a year or more, can be readily re-suspended by simple shaking or stirring."}
{"text": "The present invention relates to web page rendering, and more specifically, to optimization for rendering web pages.\nNowadays, Web 2.0 is becoming more popular, and browsers are facing more difficulties in rendering entire web pages, especially complex Web 2.0 application pages. Current rendering mechanisms of web browsers typically are not well suited for Web 2.0 single page applications (SPA), which typically include lots of Hypertext Markup Language (HTML), Cascading Style Sheets (CSS) and JavaScript together.\nTraditional web applications can cause disruption in user experience and workflow. For example, traditional web applications work by reloading the entire web page. In order to advance through a workflow, the user interacts with page elements (such as hyperlinks and form submit buttons) that cause the browser to issue a request to the server for a completely new page. Continual page redraws disrupt the user experience because the network latencies cannot be hidden from the user. There is typically a perceivable transitional jolt from one page to the next. Data of the next page is retrieved from the server, the old page is unloaded, and the new page is rendered to screen. The complete page reload that occurs on each user interaction results in unnecessary re-transmission of data over the network. This can make the overall performance of the website slower.\nA single page application (SPA), also known as single page interface (SPI), is a web application that fits on a single web page with the goal of providing a more fluid user experience akin to a desktop application. SPAs address the above issues by requiring no page reload by the browser during an application session. All user interaction and changes of the application state are handled in the context of a single Web document. The user experience becomes more continuous and fluid, and network latencies can be hidden more easily.\nIn a SPA, either all necessary code (HTML, JavaScript, and CSS) is retrieved with a single page load, or partial changes are performed which load new code on demand from the web server, usually driven by user actions. The page is not automatically reloaded during user interaction with the application, and control is not transferred to another page. Updates to the displayed page may or may not involve interaction with a server. Thus, compared to normal Web applications that still use URL to switch between various functions and pages, the single page application achieves transferring of Web application programs from server side to client side, i.e., to browsers. This mechanism brings an additional advantages of cost reduction in servers, in addition to improvements in interactivity and response speed that can be perceived by users, and lower bandwidth occupancy between servers and clients.\nHowever, since for SPA the Web application is transferred to client side, i.e., to a browser, the performance of the browser becomes important. For example, time cost on the first time rendering for a SPA web page is large, and frequent reflow and repaint via user behavior with respect to the SPA web page will bring bad performance.\nCurrently, there are several commonly used ways to improve the performance of a browser, one of which is to perform a reflow only when some reflow condition is matched, instead of performing the reflow unconditionally."}
{"text": "1. Field of the Invention\nThe present invention relates to a method for manufacturing a gallium nitride (GaN) based single crystalline substrate, and more particularly to a method for manufacturing a gallium nitride (GaN) based single crystalline substrate using a sapphire substrate.\n2. Description of the Related Art\nRecently, an optical disk field has been developed to record data at high density and to reproduce data with high resolution. In order to meet these requirements, semiconductor devices, which emit light in a short wavelength band, have been developed. A gallium nitride (GaN) based single crystalline substrate is mainly used as a material for producing these semiconductor devices which emit light in a short wavelength band. A GaN single crystal has an energy band gap of 3.39 eV, thus being suitable for emitting short wavelength blue light.\nTill now, GaN based single crystals have been grown on a substrate made of a different material using vapor growth, such as MOCVD (Metal Organic Chemical Vapor Deposition) or HVPE (Hydride Vapor Phase Epitaxy), or MBE (Molecular Beam Epitaxy). Generally, a sapphire (α-Al2O3) substrate or a SiC substrate is used as the substrate. Particularly, the sapphire substrate has a hexagonal structure the same as that of GaN and is cheaper and more stable at a high temperature than the SiC substrate, thus being widely used as the substrate for producing the above semiconductor devices.\nHowever, since a difference between lattice constants of sapphire and GaN is approximately 13% and a difference between thermal coefficients of expansion (TCEs) of sapphire and GaN is approximately −34%, strain is exerted on an interface between the sapphire substrate and the GaN single crystals, thus causing several problems such as lattice defects and cracks in the GaN single crystals. These problems cause difficulty in growing a GaN film of a high quality on the sapphire substrate, and shorten the expected lifetime of a semiconductor device produced on the GaN film.\nAccordingly, there is required a method for forming a GaN film on a GaN single crystalline layer using homo-epitaxy. However, this GaN single crystalline layer has a high vapor pressure of nitrogen, thus not being able to be used in a conventional method for producing a Si substrate or a GaAs substrate.\nTherefore, the GaN based single crystalline layer is obtained by growing a GaN bulk on a growth substrate made of sapphire or SiC using vapor growth such as MOCVD (Metal Organic Chemical Vapor Deposition) or HVPE (Hydride Vapor Phase Epitaxy), or MBE (Molecular Beam Epitaxy). For example, in case that HVPE is used, the GaN bulk can be grown on the growth substrate so that the growth thickness of the GaN bulk is several μm˜ several hundreds μm per hour. That is, the GaN bulk with a desired thickness can be grown on the growth substrate in a short period of time.\nHowever, in this case, stress is still exerted on both the sapphire substrate and the GaN layer due to the difference between TCEs of sapphire and GaN, thus causing the same problems such as lattice defects and cracks in the GaN based single crystals. In case that strain acting on the sapphire substrate by the GaN layer is less than a yielding point, the GaN layer is not cracked, but is warped toward the sapphire substrate. The warpage of the GaN layer depends on the thickness of the GaN layer. When the thickness of the GaN layer increases, the radius of curvature of the warped GaN layer is reduced and the surface of the warped GaN layer is roughened, thus causing difficulty in polishing the surface of the GaN layer.\nIn order to solve the above-described problems, there is required a freestanding GaN based substrate. The freestanding GaN based substrate is obtained by growing a GaN based single crystalline bulk on a sapphire substrate and then removing the sapphire substrate from the GaN based single crystalline bulk. Here, the sapphire substrate is removed from the GaN single crystalline bulk by mechanical polishing using diamond powder, chemical etching, or etc.\nIn case that the sapphire substrate is removed from the GaN single crystalline bulk by the mechanical polishing, stress exerted on the sapphire substrate provided with GaN single crystals grown thereon is in the range of the limit of elasticity of the sapphire substrate. Accordingly, the sapphire substrate is not cracked, but is warped. However, during the progress of the mechanical polishing, the sapphire substrate is reduced in thickness so that equilibrium of strength in the sapphire substrate is lost, thus being cracked. The cracks of the sapphire substrate may be transmitted to the GaN layer, and then the GaN layer may be cracked also. On the other hand, in case that the sapphire substrate is removed from the GaN single crystalline bulk by the chemical etching, it is difficult to obtain an etchant which has a high etching rate and selectively etches the sapphire substrate.\nRecently, there is proposed a method for growing a GaN single crystalline bulk on a sapphire substrate using HVPE and then separating the sapphire substrate from the GaN single crystalline bulk using an ultraviolet laser beam. The ultraviolet laser beam passes through the sapphire substrate with a high-energy band gap, and is absorbed by the GaN single crystalline bulk. Thus, this laser beam irradiated onto the lower surface of the sapphire substrate dissolves GaN into nitrogen gas and gallium, thereby separating the sapphire substrate from the GaN single crystalline bulk.\nThis method for separating the sapphire substrate from the GaN single crystalline bulk by means of the irradiation of the ultraviolet laser beam can be used in a small-sized substrate without causing cracks. However, in case that this method is used in a large-sized substrate with a diameter of more than 2 inches, which is generally required in manufacturing a semiconductor device, cracks form on the substrate.\nMore specifically, in case that a laser beam is irradiated on the lower surface of a growth substrate 11 made of sapphire or SiC as shown in FIG. 1, since the area of the irradiation of the laser beam is narrow (maximally 10 mm×10 mm until now), the laser beam is sequentially irradiated on limited local areas of the sapphire substrate 11 so that the laser beam can be irradiated on the entire surface of the sapphire substrate 11. Thereby, the level of stress generated by lattice mismatching and a difference between TCEs of the sapphire substrate 11 and a GaN single crystalline bulk 15 grown on the sapphire substrate 11 becomes more serious, thus causing cracks on the GaN single crystalline layer 15. Further, the cracks are propagated into the inside of the GaN single crystalline bulk 15 along cleavages, and then may cleave the GaN single crystalline bulk 15. The GaN single crystalline bulk 15 obtained by the conventional separation method is not suitably used as a substrate for manufacturing a semiconductor device thereon.\nAccordingly, there is required a method for manufacturing a GaN single crystalline substrate, in which a GaN single crystalline bulk is grown on a growth substrate and the growth substrate is separated from the GaN single crystalline bulk without damaging the GaN single crystalline bulk."}
{"text": "Cerebral blood flow (“CBF”) is among the most important physiological variables in understanding local brain metabolism. CBF provides the physiological basis for blood oxygen level dependent (BOLD) contrast, the most frequently used imaging technique for estimating changes in neural activation1 (all references cited to in this application are listed in Appendix A hereto, which is incorporated herein by reference). In addition to its fluctuations during normal brain function, CBF also changes during many of the pathological events that lead to acute or chronic brain dysfunction. Because these pathological changes typically have an earlier onset than structural biomarkers, measurement of CBF can be a useful tool in diagnosing stroke, ischemia, brain tumors, and dementia.\nThe earliest measurements of CBF used radioactive tracers.2 In fact, positron emission tomography (“PET”) measurements using radioisotopes such as O15 are still regarded as the most accurate3. However, the invasiveness and expense of radioactive tracers significantly inhibit the widespread use of PET for CBF measurement.\nDevelopments in MRI have led to an alternate method for measuring CBF, one using dynamic contrast enhancement following a rapid bolus injection of MR contrast agent. Although less invasive and expensive than radioactive tracers, this technique is still problematic, particularly in research settings. The need for venipuncture discourages some potential subjects, while others are excluded due to contraindications of the MRI contrast agent.\nArterial Spin Labeling (“ASL”) techniques replace the need for an injected contrast agent by using water protons in the plasma as an endogenous MRI tracer. ASL is both non-invasive and cost effective, and has higher spatial resolution than PET. Additionally, its clinical utility is enhanced by the fact that ASL images can be routinely acquired during the same imaging session as structural or other MRI scans and can be directly compared with the anatomical and pathological features they reveal.\nIn ASL an endogenous tracer (arterial water) is used instead of an exogenous tracer. Flowing spins are inverted (labeled) at a plane in the main arteries which is proximal to an imaging volume. The labeled spins then flow through the arterial tree into the capillary bed arriving at a particular location in the tissue with a distribution of arrival or transit times. Once in the capillary bed, the arterial water molecules then exchange, one for one, with the extravascular water molecules. Due to the accumulation of inverted spins in the extravascular component of the tissue there is a reduction in the total magnetization of the tissue. This, in turn, causes a reduction in the imaged signal in that region of the imaged plane. The reduction in signal of the tissue in the imaged plane is thus proportional to the amount of labeled spin which flowed into that plane.\nUsing ASL, in general CBF can be measured by the subtraction of two images. The first is an unlabeled or control image in which there has been no inversion of the spins. The subsequent image is labeled, i.e., the spins have been inverted in an arterial plane proximal to the imaged slice. The resultant difference between the control and labeled images is directly proportional to the flow.\nTwo common methods of conducting ASL are Continuous ASL (“CASL”) and Pulsed ASL (“PASL”). In PASL, spatially broad inversion pulses can be used to invert all of the spins in a region next to the slice of interest and the exchange between the inverted region and the imaged slice can be subsequently observed.9 In most cases PASL measures the ratio of CBF at various locations because it is difficult to precisely define the amount of inflowing spins.\nIn CASL, incoming arterial water can be continuously inverted by an adiabatic inversion pulse. Although CASL can be used to measure CBF, this process is problematic in that it relies on certain assumptions regarding the transit times and the relative T1's of blood and tissue to do so. If an image is acquired rapidly following an RF irradiation the labeled arterial water will dominate the difference signal such that any signal reduction from the water molecules that have exchanged into the tissue will not be visible. While this can be addressed by imposing a post labeling delay (PLD), for this to be successful the PLD needs to be longer than the longest transit time. This tends to make PLD's long, which results in lesser signal detection inasmuch as the difference signal continually decreases as a function of PLD length.\nIn CASL, a separate acquisition is needed to obtain a control image. Flow can then be calculated from the difference between the control image and the labeled images. Transit times between the labeling and imaging planes are not measured in CASL. Nonetheless, transit time is an important parameter in the calculation of quantitative CBF.14 To minimize the role of this unknown variable in flow quantification, post-labeling delays (PLD) between the end of labeling and the start of acquisition have been utilized, as noted. A PLD allows unlabeled spins which start flowing into the arterial tree sufficient time to wash out the labeled spins so that any labeled spins left in the arteries will not be confused with those that have exchanged into the tissue. Thus, if the PLD is longer than the transit time, the observed results are independent of transit time provided that the T1's in the blood and the tissue are the same, which is approximately true for gray matter. By varying a PLD one can investigate the effects of various transit times since an image acquired after a given PLD corresponds to integration over all longer transit times.15 \nHowever, the insertion of a PLD has several disadvantages. First, it is costly in SNR terms. Furthermore, it is vulnerable to changes in transit during activation as well as the potentially longer transit times that often occur in stroke and other cerebral vascular diseases (which can exceed the PLD). For example, a standard PLD sufficiently small not to degrade SNR beyond acceptable limits is 1 to 2 seconds, and the transit times common in stroke and other cerebral vascular diseases can be as large as 4 to 5 seconds.\nIn white matter, however, quantification of CBF using CASL encounters additional unique difficulties. The use of a PLD for measuring CBF in gray matter depends upon the assumption that the T1gray=T1blood, which, while acceptable for gray matter, is not valid for white matter. CASL also suffers from the lengthy transit time required for blood to reach white matter after being labeled in the artery. The relaxation of the labeled spins during this transit results in a much lower SNR for white matter relative to gray matter.\nIn ASL, the spins of protons in arterial water are typically inverted (labeled) at a plane in a cerebral artery. CBF then can be quantified by measuring the effect of the labeled spin when the protons enter a nearby imaging plane. However, when the transit time between the labeling and imaging planes cannot be measured, a major uncertainty is introduced into this calculation—it is difficult to discern how much of the detected signal is due to label accumulated in the tissue or label remaining in the arterial space in the tissue. While CASL attempts to reduce this uncertainty by introducing a PLD between the end of the labeling pulse and the start of data acquisition, this solution is only partially effective; there is uncertainty regarding the relative relaxation times, and it imposes an approximately three-fold cost in signal-to-noise ratio due to the longitudinal relaxation of the label during the PLD interval.\nMoreover, current ASL methods use a continuous or square waveform for spin labeling. This can be difficult for the RF amplifiers on most commercial MR scanners to generate because they are generally designed for pulsed, as opposed to continuous, applications.\nFinally, CASL uses lengthy (e.g., 1-2 second) RF pulses for spin labeling. The length of such pulses tends to place strain on the RF transmitters in most MRI scanners that is typically relieved through the use of transmit-receive head coils. However, receive-only coils offer greater sensitivity and are more widely available making this solution non-optimal. Additionally, long RF pulses introduce confounds from magnetization transfer (MT) effects whereby the spins in the imaged plane are saturated by the long RF pulse. To address this, CASL requires that the acquisition of a control image be done in the presence of a RF pulse which mimics the MT effects of labeled images but does not invert the arterial spins. However, while this approach eliminates MT effects, it introduces additional sources of error, particularly motion artifacts.\nThus, what is needed in the art is a method of blood flow and perfusion measurement that can overcome the problems of conventional methods, and that can allow for the direct measurement of transit time distribution while acquiring ASL data.\nWhat is further needed in the art is a method of blood flow and perfusion measurement that can add certainty to CBF quantification and obviate the need for PLD's control image acquisition or reliance on less sensitive signal receiving hardware."}
{"text": "1. Field of the Invention\nThis invention relates to improved apparatus for error correction of digital data, and in particular to correction of digital data representing a pictorial image.\n2. Description Relative to the Prior Art\nWhile much digital data, such as that derived from computation and general data processing has no redundancy, data generated by digital processing of images is characterized by substantial amounts of redundant information. The established method of converting images into digital data entails projecting the image onto a sensor having an analog output, and then scanning the sensor and digitizing its output. The redundancy inherent in the image is carried over into the resultant digital data stream; the digital values themselves exhibiting the redundancy of the original pictorial image.\nReferring to FIG. 1, a scene 10 is scanned in horizontal segments by an optical scanning device 12 whose output is a raster 14 containing the pictorial information in the form of electrical analog signals. The scanning device 12 may be any one of the optical scanners known in the art; for example, the laser scanner, kinescope scanner, or CCD sensor and associated scanning circuits. The electrical signals are then sequentially digitized by means of a sampling analog to digital converter 16 whose output is the digital representation of the image intensity on a line by line basis. Such image samples are conventionally referred to as \"pixels\". Generally, the digitizer 16 output comprises the pixel value expressed in binary form wherein eight bits per sample correspond to a 256 level gray scale of image intensity. A digital eight bit sample is usually designated as a \"byte\", however, a byte may be defined to contain any discrete number of bits. This binary coded representation of the digital values of the samples is called \"pulse code modulation\", or PCM. Usually the image is digitized to allow either electromagnetic transmission or magnetic recording of the PCM encoded samples.\nIn FIG. 1, the line 18 represents a left-to-right scan across the \"sky\" portion 19 of the scene 10, and it will be appreciated that the image intensity across such a scan will be essentially constant. Similarly, the intensity of the scan of the line 20 across the sky remains essentially constant until the edge of the \"cloud\" 22 is encountered, and then the image intensity changes to another essentially constant value for the duration of the scan across the \"cloud\". Upon completion of the scan of the cloud, the intensity and the sample values then revert to the values representative of the \"sky\" 19. Such scans result in sample values that are highly redundant because there is little detail in the portions of the scene being sampled. On the other hand, the detailed portions of the scene which consist of rapidly changing intensities result in samples of limited redundancy; during the scan 23 across the \"tree\" 25, rapidly varying sample values of intensity are generated. Thus, the data is characterized by non-redundant bytes derived from scans of detailed portions of the image such as the \"tree\" 25, and by redundant byte values derived from scans of the \"sky\" 19 or the \"cloud\" 22 where the intensity changes little, if at all.\nThe present inventor's copending application, Ser. No. 023,327, U.S. Pat. No. 4,761,782 focuses attention on the use of image redundancy for error correction after the transmission or storage of the digitized image data. Aberrations such as fading in an electro-magnetic transmission channel, or dropouts in a magnetic tape storage channel cause signal loss and upon recovery of the data, attendant errors. These errors result in easily discernible, detrimental streaking in the recovered image. The teaching of the above copending application ameliorates this problem by utilizing the available image redundancy as the basis for error correction in the recovered image data. It teaches the grouping of the digitized bytes representative of an image into data blocks whereby the redundant information present in the image becomes replicated in the data block and results in some of the bytes having essentially the same, or nearly the same, digital values as other bytes within the same block. Prior to transmission or storage, each block is analyzed for redundant sample values, i.e. values which are either identical or which only vary within predetermined boundaries, and check bits are accordingly affixed to the block. The configuration of the check bits defines the nature and the extent of the redundancy in the bytes comprising the block. The block containing data and check bits is then transmitted or stored. The recovered data, which has been subject to induced errors during transmission or storage, is analyzed to ascertain whether the bytes comprising the block still contain the same redundancy specified by the affixed check bits. If an error has occurred, the resultant redundancy of the received data will generally be reduced; at least one of the recovered pixels in the block will be different from what it originally was due to the error. The above mentioned copending application teaches re-encoding the received data to derive new check bits based on the received data which are then compared to the original check bits as received. This comparison generates a \"syndrome\" having the value \"0\" if the check bits derived from the received data are the same as the original check bits, and having the value \"1\" if the re-encoded check bits do not agree with the original check bits. A syndrome of value 1 flags an error and the fact that the data in the block has changed during transmission. As described in the copending application, if an error has occurred, the received data bytes are then compared among themselves to ascertain which data pattern is characteristic of the majority of received bytes. This majority pattern is assumed to be the correct pattern for all the bytes of the block including the bytes in error. The bytes in error are then corrected in accordance with the instructions inherent in the configuration of the original check bits. The correction restores the data patterns to their original configuration, and, attendantly, restores the amount of redundancy that was originally present in the block.\nThe check bits are not derived from the actual values of the data bits as is characteristic of many codes known in the art. The check bits reflect the fact that \"patterns\" of data bits within the block are the same, but the check bits do not explicitly define the structure of these patterns.\nIn considering the teaching of the prior art, it is advantageous to provide a listing of various check bit symbols used in describing the technique.\nSx=check bit x derived from data before transmission, PA1 SxR=received value of check bit Sx after transmission PA1 Sx*=re-encoded value of Sx derived from data as received after transmission PA1 SCx=value of syndrome calculated from SxR and Sx*, where for m check bits PA1 x=0,1,2 . . . m-1. PA1 127=01111111 PA1 128=10000000 PA1 132=10000100 PA1 129=10000001 PA1 127=01111111.\nA prior art embodiment employs a two bit check pattern, S1, S0, which is attached to each block of n pixels to provide an 8n+2 bit block. Check bits are assigned in accordance with the following rules:\nTABLE I ______________________________________ S1 S0 CONDITION ______________________________________ 1 1 Three most significant bits of all n pixels form indentical patterns 1 0 Only the two most significant bits of all n pixels form identical patterns 0 1 Only the most significant bits of all n pixels are the same 0 0 None of the above ______________________________________\nThe above assignment of check bits reflects the degree of redundancy carried by the three most significant bits of the n pixels comprising a block. For example, the check bits 1,1 indicate that the n pixels have the highest redundancy in their three most significant bits, i.e. the three most significant bit patterns for all n pixels are the same. On the other hand, the check bits 0,0 indicate that there is no redundancy among the most significant bits of the n pixels, i.e. all of the patterns are different. The intermediate cases of 1,0 and 0,1 reflect in-between conditions of redundancy.\nIt will be appreciated that the prior art technique depends upon the image redundancy as expressed in the three most significant binary bits of the pixel binary representation. Because of the structure of binary representation, it is possible that certain values of redundancy will fail to encode in a manner correctly representing the redundancy value. For example, a block may consist of digitized image values of 127, 128, 132, 129, 127. Encoded into an eight bit binary representation, these values are:\nWhile it is clear by consideration of the actual magnitudes of these pixels that the information is highly redundant, neither the three most significant bits of the binary values of the five pixels are the same, nor are the two most significant bits, nor are the most significant bits themselves, the same. Due solely to the structure of binary representation, these five pixels appear to be completely non-redundant rather than highly redundant. This occurs at the transition between the numbers 127 and 128 in binary representation because all the bits switch in value with the result that values just above the transition point and values just below it have no similarity in the high order bit positions. A similar \"boundary\" problem also occurs for pixels having digitized binary values in the region of 256, and 64. It will be appreciated that because of this characteristic large undetected errors may occur at these boundary points. For example, it is clear that if the most significant bit of the binary number 128 is subject to error and changes from 1 to 0, the resultant pixel value instead of being 128 is 0; a level change of 50 percent of full scale, (i.e. 256), that is not corrected in the prior art due to coding failure at the boundary.\nAn analysis of typical images containing redundancy shows that from 7 to 15 percent of blocks may fail to incorporate a block's true redundancy value when utilizing the techniques of the prior art. This is partially due to the redundancy characteristics of the image and partially due to the above described boundary problem. The mere occurrence of pixel values which straddle a boundary region is not sufficient in itself to cause errors in the reproduced image. It is necessary that a transmission error or, say, a magnetic tape storage dropout error simultaneously occur with the incidence of the pixel value in the boundary region. Because the joint probability of these two conditions occurring simultaneously is very small, acceptable error correction in many applications is provided by the technique of the prior art. For the processing of the highest quality images wherein the occurrence of any large pixel errors, even if rare, is unacceptable, additional correcting power is required. The present invention provides that power by completely eliminating the source of the errors due to the boundary problem."}
{"text": "1. Field of the Invention\nThe present invention relates to electrical connectors and more particularly to modular connectors for use in connecting a daughter printed wiring board to a mother printed wiring board.\n2. Brief Description of Prior Developments\nIn the manufacture of computers and various other electronic assemblies, daughter boards are commonly connected to mother boards by means of a connector having a receptacle having a plastic housing and a first and second face wherein terminals are connected in one face to the daughter board and at the other to a header connected to the mother board. Various arrangements have been suggested to ground such connectors to the mother or daughter boards but such arrangements have tended to complicate the construction of the connector. A need, therefore, exists for simple and inexpensive means for grounding connectors between mother and daughter boards. There is also a need for such a connector which reduces crosstalk and increases band width.\nThe shield arrangement of the present invention comprises an external conductive shield and an internal conductive shield. The external shield encloses the electrical connector and has at least one inwardly directed projection. The internal shield at least partially resides within the electrical connector and has at least one feature engaging the at least one projection.\nThe electrical connector of the present invention comprises: a housing; a plurality of contacts; an external conductive shield; and an internal conductive shield. The contacts extend through the housing. The external shield encloses the housing and has at least one inwardly directed projection. The internal shield at least partially resides within the housing and has at least one feature engaging the at least one projection.\nThe receptacle electrical connector of the present invention comprises: a housing; a plurality of contacts; an external conductive shield; and an internal conductive shield. The housing has a front face with a plurality of apertures arranged in a plurality of columns. The contacts are also arranged in a plurality of columns in the housing, each contact having a mating end adjacent a corresponding openings. The external conductive shield encloses the housing and has at least one inwardly directed projection. The internal conductive shield at least partially resides within the housing at a location between adjacent columns of contacts, and has at least one feature engaging the projection."}
{"text": "This invention relates to innerseals for containers, and, more particularly, innerseals that are suitable for containers made of glass.\nHeat-sealable innerseals have been found to be useful in sealing containers for liquid products, such as, for example, motor oil, brake fluid, antifreeze, household ammonia, liquid detergents, etc., which products present technical problems with respect to leakage. Heat-sealable innerseals have also become popular for providing the tamper-resistant innerseals that are required by the Food and Drug Administration for over-the-counter drugs.\nHeat-sealable innerseals generally comprise a layer of aluminum foil bearing a coating of heat-sealable material thereon. The innerseals are inserted into caps, and the resulting assemblies are then supplied to the packager. The cap is then placed by the packager onto the filled container, with the coating of heat-sealable material being in contact with the lip or rim of the container. The container then passes under an induction heater which generates heat through the aluminum foil, thereby melting the heat-sealable coating, and causing the innerseal to bond to the container.\nHeat-sealable innerseals applied by induction heating have been found to be especially useful with plastic caps and plastic container systems, as induction heating does not heat the plastic material excessively. Heat-sealable innerseals can also be used with metal caps in an induction heating system where the metal cap itself conducts heat to melt the heat-sealable material.\nIn the case of glass containers, however, heat-sealable innerseals will lose adhesion within days, even hours, however, upon exposure to moisture. This loss of adhesion occurs because the sodium ion, which is an essential component of many popular glass containers, gradually leaches out of the bulk of the glass and reacts with water in the air or in the container to form a layer of sodium hydroxide on the surface of the glass. It is this layer of sodium hydroxide that reduces the adhesion between the heat-sealed coating of the innerseal and the glass surface of the container.\nOne innerseal material in current use that overcomes the aforementioned problem is an ionomer sold under the trademark \"Surlyn\". However, this ionomer has a disagreeable odor and, at high temperatures, portions thereof leach into the liquid contents of the container."}
{"text": "1. Field of the Invention\nThis invention relates to an improved liquid detergent having added softener and brightener for use in laundering fabrics. In particular, it relates to a storage stable, single-phase laundry detergent-softener-brightener composition free from phosphates and possessing enhanced softening and brightening properties.\n2. Description of the Prior Art\nIn the recent past, laundry was washed by simply placing it in a machine filled with detergent, running the machine through a wash cycle and then through a rinse cycle.\nLater, cationic fabric softeners were added to the rinse cycle to improve the feel of the laundry and reduce the static charges built up on fabrics during drying. These cationic fabric softeners were not blended with liquid detergents, since they became unstable in the presence of the detergents and tended to form distinct phases during storage. The softeners were not successfully employed in the wash cycle, since they tended to deposit on the clothing as a yellow film. This film also tended to coat fabrics with a water-repellent finish after repeated use. In fabrics, particularly those which needed to breathe or absorb water, this was a serious defect. Fine fabrics were recognized to be particularly susceptible to yellowing, clinging and filming from detergent washing with softening treatment.\nIn an attempt to correct these defects, anionic detergents and phosphate builders were added to detergents to complex the softeners and prevent their yellowing and filming of fabrics. However, these attempts proved unsuccessful, since the softening effect of the fabric softeners, notably the quaternary ammonium salts of \"quats\", was seriously reduced. In addition, it has been recognized, recently, that phosphates can adversely effect the ecology. Consequently, there has been much clamor for detergents, which are free from phosphates.\nModern fabrics, including nylons, polyesters, cotton and blends thereof and particularly fine fabrics, are expected to be bright in appearance after washing. Commonly used detergent fluorescent brighteners, such as DASC brighteners (DASC-ASTM designation--those derived from diaminostilbene disulfonic acid-cyanuric chloride) have two sulfonic acid groups per molecule. These brighteners have been suggested for use with detergents and softeners. However the disulfonate brighteners tend to form unstable opaque products with very poor brightening properties, when combined with laundry detergents and cationic softeners. Monosulfonate, NTS-type stilbenesulfonic acid brighteners have been proposed, but provide very poor brightening when employed with detergents and softeners. Diaminostilbene hexa-sulfonate brighteners are used in coatings in the paper industry, but are not recommended for use as brighteners with standard detergents owing to their high level of water solubility. Such hexa-sulfonate brighteners and cationic softeners. Monosulfonate, NTS-type stilbenesulfonic acid brighteners have been proposed, but provide very poor brightening when employed with detergents include those disclosed in British Pat. No. 941,993 which utilize detergents and DASC brighteners, but no fabric softeners. These compositions are merely powder mixtures, not liquid compositions.\nIn U.S. Pat. No. 3,625,891, a blend of three different quaternary ammonium (\"quats\") fabric softeners were employed to overcome detergent compatibility problems. The softener and detergent were added separately to the wash.\nIn U.S. Pat. No. 3,644,203 a minor amount of anionic surfactant was employed with a cationic fabric softener together with a nonionic surfactant and a conventional phosphate builder, forming a multifunctional detergent, in which the anionic surfactant served to neutralize the cationic softener. Optical brighteners were optionally disclosed, but not those employed in this invention.\nA nonphosphate softener-detergent composition is set forth in U.S. Pat. No. 3,959,157. A DASC brightener is employed having only two sulfonate groups. To stabilize the \"quat\" softeners, up to 20% of an alcohol is disclosed to be necessary. This adds to the expense of the composition. Further, the level of \"quat\" softeners employed in the patent is from about 3% to 15% by weight of the composition. It is most desirable to reduce this level, while maintaining softening properties, to reduce the costs of the product to the consumer and reduce yellowing and filming tendencies."}
{"text": "Full well capacity is one of the major characteristics of Complementary Metal-Oxide-Semiconductor (CMOS) image sensors. The full well capacity is a measurement of how much charges an image sensor may store before the charges overflow. The full well capacity determines the dynamic range of the image sensors. A high full well capacity means that the respective image sensor may have a great difference between the brightest level and darkest level of sensed signals. Hence, the full well capacity needs to be known to determine the quality of the image sensor chip. If necessary, the manufacturing process of the image sensor chips may be adjusted when the respective full well capacity does not meet design requirement."}
{"text": "In a camshaft free combustion engine a pressure fluid, such as a liquid or a gas, is used to achieve a displacement/opening of one or more engine valves. This means that the camshafts, and related equipment, that conventional combustion engines use to open engine valves to let air in respective let exhaust fumes out from the combustion chamber, has been replaced by a less volume demanding and more controllable system.\nIn an engine that is constructed for significant angular momentum outputs, the pressure in the combustion chamber is increasing proportional to an increased angular momentum output, and the force that is required to open the valve actuator to open the, in relation to the combustion chamber inward opening, engine valve consequently also increases proportional to an increased angular momentum output. At high numbers of revolutions, such as 6-8000 rpm, a very fast opening of the engine valve is also required for the filling of air respective evacuation of exhaust fumes from the engine cylinder not to be restricted. These requirements, i.e. the need for an extremely fast opening at high frequencies in a high performance engine having high counter pressure in the combustion chamber of the engine at the opening of the exhaust valves, require the pressure of the pressure fluid upstream of the valve actuator to be high, in the order of 8-30 bar. Downstream the valve actuator, the pressure fluid has a lower pressure, in the order of 3-6 bar.\nAt high numbers of revolution and high engine loads, the pressure difference between the low pressure side and the high pressure side should be in the order of 15-20 bar to achieve a correct operation of the valve actuators, and when the engine is idle running, or at low numbers of revolution and low loads, the pressure difference between the low pressure side and the high pressure side only needs to be in the order of 2-5 bar. The lower pressure difference at low numbers of revolution is desirable when the pressure is increased by way of a compressor from the low pressure side to the high pressure side, and at the pressure increase, energy consumption occurs that increases concurrently with increased pressure on the high pressure side.\nIn situations that require fast acceleration and/or very fast change from low numbers of revolution and low load to high numbers of revolutions and high load, for example when entering on a busy main road or at a sudden overtaking of a slow moving vehicle, the pressure difference between the low pressure side and the high pressure side must immediately be increased to achieve the acceleration that the driver requires. A conventional compressor is dimensioned to be able to generate pressure differences with greatly varying magnitude, they are however not dimensioned to satisfy the requirement of immediate shifts between separate great pressure difference levels and pressure fluid flows.\nFurthermore, there is inertness in the present systems to go from a great pressure difference to a small pressure difference, i.e. when the vehicle again is operated at low numbers of revolution after the short/temporary rise in numbers of revolutions/engine load, it will take time before the pressure difference and thereby the high energy consumption has decreased to a desired level."}
{"text": "1. Field of the Invention:\nThe present invention relates to an industrial robot of the type wherein a wrist section thereof has a freedom of movements along three axes.\n2. Description of the Prior Art:\nIn a known industrial robot having a freedom of movements along six axes, a twist body is rotatably carried at one end of an arm body, a bend body is carried at one end of the twist body for rotation about an axis which is slanted with respect to the rotational axis of the twist body, and a swivel body is carried on the bend body for rotation about an axis which is slanted with respect to the rotational axis of the bend body. The known robot is accompanied by a problem that the working range of the bend body is diminished to avoid the interference between the twist body and the bend body.\nIn order to enlarge the working range of the bend body, it is effective to provide the front end of the twist body with an eccentric cylindrical portion whose axis is eccentric a predetermined amount from the rotational axis of the twist body and to coaxially and rotatably carry two rotational shafts which transmit rotational power respectively to the bend body and the swivel body. If this configuration is used, however, it may become practically difficult to assemble these long rotational shafts into the twist body, with the rotational shafts being held in coaxial relation. Particularly, it may not be easy and may result in taking a long time to assemble the rotational shafts directly into the eccentric cylindrical portion of the twist body.\nFurther, if the long rotational shafts for the bend and swivel bodies were disposed at a position which is eccentric from the rotational axis of the twist body, there may be raised a problem that the moment of inertia required to rotate the twist body has to be increased."}
{"text": "The present invention relates to roadside beacon systems, and more particularly to a roadside beacon system used for calibrating a vehicle position in a navigation system in which, after data representing the start point of the vehicle are inputted, the current position of the vehicle is calculated by measured vehicle speed and direction data to display its present position and also used for carrying out data communications between the vehicle and the roadside antenna.\nA vehicular navigation system has been proposed in the art in which a small computer and a small display unit are installed in a vehicle; a road map stored in a memory device such as a compact disk is displayed on the display unit; speed data outputted by speed sensors and direction data outputted by direction sensors are utilized to calculate the present position of the vehicle and to determine the present direction of movement of the vehicle at all times, and according to the results of the calculations, a symbol representing the vehicle is indicated on the road map displayed on the display unit.\nWith such a navigation system, the operator of the vehicle can visually detect the present position and direction of movement of the vehicle with ease and thus can find a destination readily.\nHowever, the above-described navigation system suffers from a potential for inaccuracies which are due to errors accumulated in speed sensors and direction sensors as the vehicle travels a long distance. Therefore, when the distance covered by the vehicle exceeds a certain value (which is not always a constant value and is dependent on the amount of errors accumulated in the speed sensors and the direction sensors and also environmental conditions), the vehicle position displayed on the display unit may be greatly shifted from a true position; that is, the navigation system becomes unreliable and the driver may lose his way.\nIn order to eliminate the errors in navigation systems, a so-called \"roadside beacon system\" has been proposed. In this system, roadside antennas are installed along roads at intervals shorter than those with which the errors thus accumulated reach serious values. Each of the antennas radiates signals including position data of antenna location and road direction data to cover a relatively small area. The signals thus transmitted are received by a mobile antenna installed on the vehicle and applied to the navigation equipment whereby, according to the signal thus received, the position and the moving direction of the vehicle are calibrated.\nIn the roadside beacon system described above, the accumulation of errors is attempted to be maintained at less than aggravating value, whereby an exact present position of the vehicle can always be displayed. Accordingly, the navigation system is reliable. Furthermore, the roadside beacon system has the advantage that, when a roadside antenna is installed, for instance, near a railroad or railroad crossing where the direction sensors are liable to produce a large error, errors attributed to external factors can be corrected effectively.\nIn the above-described roadside beacon system, however, the roadside antennas are of considerably high directivity, transmitting signals including position data and road direction data at all times, and the vehicle receives the beacon signals only when it passes through a covered area by the signals. If one attempts to cover large areas with each beacon antenna in order to get large communication area it becomes difficult to accurately detect the vehicle's position, because of wide spread field distribution which results in difficulties in exactly distinguishing the antenna position.\nThis will be described in more detail. The fundamental function of the roadside beacon system is to transmit signals including the position and road direction data to the navigation system in the vehicle. It may also be required for the roadside beacon system to further transmit data to perform the following functions in practical use:\n(1) transmission of traffic information as to traffic congestion, construction work on road, detours, etc., in the area around the roadside antenna;\n(2) information for a detailed map including the position of buildings or houses with names in the area around the roadside antenna;\n(3) information for a road map of a relatively large area including the position of the roadside antenna to renew the road map on the display unit, and\n(4) two-way communications between the roadside antenna and the vehicle.\nAbove mentioned applications can be fulfilled only with using an enlarged communication area covered by the roadside antenna in order to realize large volume data transmissions.\nHowever, if the communication area covered by the roadside antenna is enlarged, the detection of the vehicle position which is the primary object of the navigation system, becomes inaccurate.\nIn addition to the spread field distribution in the enlarged communication area system, there are environmental obstacles to electro-magnetic waves such as vehicles moving nearby, buildings along a road and so on, causing the signal received by vehicle to fluctuate largely. An example of measured field distribution fluctuation caused by the above mentioned factor is shown in FIG. 14a. The fluctuation of signals is the results of receiving waves scattered by or reflected from such obstacles being different in amplitude and phase from one another. In other words, multi-path fading occurs in the signal. Accordingly, calibration of the vehicle position using the signal thus received may involve an unexpected error. That is, the signal received by the mobile antenna may have a high level at a position actually farther away from the roadside antenna, resulting in large positional errors.\nThe present applicant has filed a Japanese patent application for a roadside beacon system in which, to eliminate this difficulty, \"split-beam\" antennas having a radiation pattern whereby the electric field strength abruptly decreases directly in front of the antenna are employed for the roadside antennas. When the abrupt decrease of the received signal strength is detected, it is determined that the vehicle is directly confronting the roadside antenna, and the vehicle position data and the moving direction data are calibrated according to the data received just before the detection of the abrupt decrease point. Hence, the effects of multi-path fading on the received signal strength distribution are eliminated.\nIn this connection, a Japanese patent application has also been filed for a method in which roadside antennas are installed at elevated positions above the road, and the major lobe of the mobile antenna is extended upwardly. Also, a Japanese patent application has been filed for a method in which a roadside antenna is connected to the lower portion of a structure which extends over the road, and the major lobe of the mobile antenna is extended upwardly.\nIn these roadside beacon systems, a variety of data can be transmitted between the roadside antenna and the vehicle over a wide range with no difficulty, and the accuracy of detection of the vehicle position with respect to the roadside antenna can be considerably improved.\nHowever, these roadside beacon systems may suffer from the following difficulty: In the case where large obstacles like truck and buses which are considerably larger than the vehicle installed with navigation system, are existing around the roadside antenna and the vehicle is in the communication area as shown in FIG. 2, or in the case where the roadside antenna is installed under the bridge structure which is constructed over the road illuminating the road from upside and a large obstructing vehicle or vehicles are moving nearby the vehicle installed with the navigation system as shown in FIG. 3, the large obstructing vehicle will shield or scatter the radio waves resulting in the deterioration of positioning accuracy of the vehicle.\nExperimental results will be described in more detail. In the experiment the large obstructing vehicle traveling the central lane shown in FIG. 2 is 8.12 m long, 2.2 m wide and 3.5 m high, the vehicle with navigation system is 1.0 m high, the roadside antenna is 5.0 m in height, and the inclination angle of the main radiation beam of the roadside antenna is 30.degree.. In this case, the fluctuation in the signal level received by the mobile antenna is as shown in FIG. 4A. As is apparent from the comparison with the fluctuation of the signal level (as shown in FIG. 4B) in the case of clear environmental condition in which no substantial electromagnetic obstacles such as large trucks, buses and so on, the signal level variation (FIG. 4A) shows characteristics that an abrupt level decrease attributed to the shielding of the radio waves by the large obstacles and ripple components of up to 10 dB are superposed on the broader and smooth signal level distribution curve determined by the directivities of the roadside antenna and the mobile antenna. Therefore, the accuracy of detection of the vehicle position is somewhat deteriorated.\nAnother experimental data will now be considered where, in FIG. 3, a large vehicles is traveling along a center lane just below a roadside antenna which is 6 m in height, and the vehicle having the navigation system is moving along an adjacent lane. In this case, the variation in level of the signal received by the mobile antenna is as shown in FIG. 4C. As is apparent from a comparison with the variations the signal level received in the case no large vehicle is moving parallel to the vehicle with the navigation system (as shown in FIG. 4D), the signal level variation of FIG. 4C is such that ripple components of the order of 3 dB attributed to multi-path fading are superposed on the broader and smoother signal strength distribution curve determined by the directivities of the roadside antenna and the mobile antenna.\nIn both of the above-described experiments, the frequency of waves which carry transmitting data is set at about 2.5 GHz.\nThus, in both of the above-described cases, the signal received by the mobile antenna includes ripple components of the order of 3 to 10 dB. However, this causes no problem in data transmission; that is, the data can still be transmitted to the mobile antenna at high speed, because of enough C/N characteristics obtained at the bottom of signal fluctuations.\nOn the other hand, in the case where a split beam signal from the roadside antenna is used to detect the vehicle position, the following difficulties are involved: Since the split beam signal radiated from the roadside antenna sharply decreases just in front of the roadside antenna and at positions sufficiently far away from the roadside antenna, it is insufficient for detection of the vehicle position merely to detect the decrease in level of the received signal. That is, only sharp fall of the signal level which occurs when the vehicle goes just in front of the roadside antenna.\nTo satisfy the above requirement, a navigation device has been proposed which has been designed as follows: In the device, the received signal is applied to a low-pass filter so that abrupt variations are removed; i.e., the signal is converted into a signal having a more gently curved characteristic. The signal thus obtained is divided into two signals. From one of the two signals, it is detected that the vehicle is nearing the roadside antenna. When it has been detected that the vehicle has arrived at the region adjacent the roadside antenna, the other signal is applied through a gate circuit to a position detecting section so that the point where abrupt falling of the received signal level occurs can be detected. Thus, the navigator device can detect an abrupt decrease of the received signal level which corresponds to the position just in front of the roadside antenna. However, the device is still not sufficient in that, since the abrupt variations in the level of the received signal are removed by the low- pass filter, when the vehicle is driven at a high speed, the dip in the level of the received signal is not so deep, and therefore the position detecting section may not be able to detect the abrupt decrease in level of the received signal.\nThe above-described difficulty may be eliminated by removal of the low-pass filter. However, for reasons clear from the above discussion, doing so results in the difficultly that the received signal level decreases not only at the aforementioned abrupt decrease point P (cf. FIG. 4E) where the level should decrease, but also at other points P.sub.1, P.sub.2 and P.sub.3. Therefore, the position detection becomes inferior.\nMoreover, the level of the signal received by the mobile antenna varies with the lane on which the vehicle moves; that is, it changes with the distance between the roadside antenna and the mobile antenna. Furthermore, the level depends on the inherent fluctuations in sensitivity of the roadside device and also the mobile device. Therefore, it is very difficult to set the fixed reference signal level at which the positioning of vehicle can be determined by detecting the falling of signal level to the reference point.\nExperimental results will be described in more detail. For instance in the case where, on a three-lane highway, a split beam antenna is installed as a roadside antenna alongside the first lane, and an antenna of which main radiation lobe extends upwardly is used as the mobile antenna on a vehicle, the signal received by the mobile antenna on the vehicle traveling in the first lane is the highest in level (A in FIG. 5 and FIG. 6A), the signal received by the mobile antenna of the vehicle traveling in the third lane is lowest, and that received by the mobile antenna of the vehicle in the second lane falls between those two signals in level (B in FIG. 5 and FIG. 6B).\nIt is observed that there exists about 10 dB variations in receiving signal level according to the lane the vehicle is moving; however the sharp falling of the signal level just in front of the roadside antenna can be seen on any lanes.\nThis difference was observed in the experiment with only one set of roadside devices and the same mobile devices respectively. Considering a practical case in which a lot of beacon antennas are employed and a huge number of vehicles are involved, the variations in the signal level may actually become larger than 10 dB. Therefore, it is substantially impossible to set an absolute reference level for detecting the abrupt decrease the received signal."}
{"text": "Solar cells convert light into electricity based on photoelectric effect and are used as for providing clean energy without pollution. Nowadays, solar cells have been widely researched and manufactured and are gradually adopted to replace traditional electric power.\nGenerally, solar cells are formed with a semiconductor device, such as monocrystalline silicon, multicrystalline silicon, amorphous silicon or other thin film semiconductor devices which include a p-n junction. The method of forming the p-n junction includes doping n-type or p-type ions into semiconductor device by diffusion or ion implantation. When the semiconductor device is illuminated by light, electron hole pairs are generated in semiconductors, such as silicon, and electrons and holes are separated drifted to form a current by an internal field formed by the p-n junction. In order to collect the current generated by the solar cells, metal electrodes are formed on surfaces of the solar cells. As shown in FIG. 1, the metal electrodes formed on the solar cells 10, 20 include back electrodes 12, 22, finger electrodes 13, 23 which are formed on the front side 11, 21, and busbar (or buss bar) electrodes 14, 24 which are formed on the two terminals of the finger electrodes 13,23.\nThe finger electrodes 13, 23 and the busbar electrodes 14, 24 are usually formed by coating silver paste on surfaces of the solar cells 10, 20 by screen printing. Besides, the finger electrodes 13, 23 are usually made with smaller widths and arranged in parallel with a predetermined interval to minimize shading effect. The back electrodes 12, 22 also can be formed by screen printing and is usually made of aluminum. However, the back electrodes 12, 22 are different from the finger electrodes 13, 23 that they can be formed by coating aluminum on a whole surface of the solar cells 10, 20.\nA single solar cell usually provides an output voltage of about 0.5 volt to about 0.7 volt, and therefore a plurality of solar cells are connected in series to constitute a solar module to provide a required operating voltage. As shown in FIG. 1, neighboring solar cells are usually series connected with their front and back electrodes via ribbons 30. For example, a terminal of the ribbon 30 is connected to the busbar electrode 14 formed on the front side of the solar cell 10 while the other terminal of the ribbon 30 is connected to the back electrode 22 of the solar cell 20 and therefore a solar module 1 is formed with serially-connected solar cells, hereinafter a “series-connected type solar module.”\nThe above process requires steps of forming metal electrodes and soldering ribbons. Since traditional soldering process is usually performed at a higher temperature, micro cracks may be generated on the surface of the crystalline silicon solar cells because of the thermal expansion/shrinkage coefficient difference between metal electrodes and crystalline silicon substrate."}
{"text": "The invention relates generally to the field of vehicle propulsion systems. In particular, the invention concerns an on board recharging system for a battery powered electric vehicle.\nAs concerns surrounding traditional energy sources persist, investigation into alternative forms of energy are becoming increasingly important. In particular, environmental and political concerns associated with combustion based energy systems can not be ignored. In an effort to reduce dependence on these types of energy supplies and methods, therefore, the focus is likely to be on devices capable of generating electricity by consuming plentiful or renewable supplies of fuels without requiring combustion.\nIn addition to being utilized for the generation of electricity, however, combustion based fuels are most commonly utilized for driving motor vehicles. Indeed, motor vehicles are among the chief consumers of combustion based fuels and, therefore, they are among the chief contributors to pollution problems associated with such fuels. So, while alternative energy sources such as nuclear and hydroelectric systems may be suitable for large scale electric utilities, for obvious reasons they do not present an ideal solution to the problems associated with motor vehicle power requirements.\nAlternatives to internal combustion engine powered motor vehicles have included various types of battery powered electric vehicles. A problem associated with most known battery powered vehicles, however, is the constant need for recharging which requires taking the vehicle out of service.\nFuel cells, therefore, have been explored as a means for powering electric vehicles and reducing the constant need to recharge the vehicle from off-board sources. Fuel cells electrochemically convert hydrocarbons and hydrogen to electricity- No combustion reaction is needed. Because they are compact in structure, high in efficiency, and low in pollutants, fuel cells provide several advantages for use in vehicle transportation. Accordingly, they appear to be well suited for adaption for motor vehicle applications.\nA drawback associated with known fuel cell systems, however, is that they are not economically viable for applications in which the power rating of the fuel cell must meet propulsion demands. In motor vehicle applications, for example, a fuel cell system designed to provide sufficient power required by the vehicle for cruising, let alone for peak surge, would be prohibitively expensive. While various known systems have attempted to exploit the advantages of designating a surge battery to meet peak demand in motor vehicle applications, none has satisfactorily overcome the economic problems.\nIt is an object of the invention, therefore, to provide a system for powering a motor vehicle which does not require a combustion reaction. It is another object of the invention to provide such a system which has a range comparable with that of traditional combustion engines without requiring interruptive recharging from off-board sources. It is still another object of the invention to provide an economically feasible system for powering a motor vehicle which can accommodate typical motor vehicle surge and range demands."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of manufacturing an electronic device having a fine pattern such as a semiconductor device or a wiring substrate device.\n2. Description of Related Art\nIn fabricating a semiconductor circuit, there are repeatedly employed a film forming step such as CVD (Chemical Vapor Deposition), an impurity layer forming step such as ion implantation, a lithography step for forming a resist pattern, and an etching step. Microsizing a circuit pattern and improving the dimensional accuracy of the circuit pattern are effective as a method for improving the operating speed of a semiconductor circuit or improving the degree of device integration. Recently, both microsizing and improvement of dimensional accuracy have been promoted actively. The microsizing of a pattern mainly depends on lithography, so lithography now occupies an extremely important position in the manufacture of a semiconductor device.\nIn the lithography technique, a projection exposure system is primarily used. A pattern of a photomask mounted to the projection exposure system is transferred onto a semiconductor wafer to form a device pattern.\nTo meet the recent demand for higher device integration and a higher device operation speed, the microsizing of the pattern has been promoted. With this as background, an exposure method called a halftone phase shift method is employed. The halftone phase shift method indicates a mask wherein a film (called a halftone film) translucent to exposure light is formed on a transparent substrate (blank). The transmittance through the film of exposure light is adjusted usually to a value in the range from 1% to 25%. Exposure light for passing through the film is adjusted so as to be phase-inverted in the absence of the film. A phase difference which brings about the highest resolving performance is 180° or an odd-number multiple thereof. But insofar as the phase difference is inside of 90° above and below 180°, there is obtained a resolution improving effect. It is generally known that the use of the halftone phase shift mask brings about an improvement of approximately 5% to 20% in resolution. Particularly, the higher the transmittance of the halftone film, the more outstanding the phase enhancing effect, i.e., the higher the resolution becomes. On the other hand, an image called a sub peak is apt to be formed in a field portion which should be a light shielding portion, and causes a transfer defect. This is a phenomenon caused by simultaneous enhancement of the phase of interference light from an adjacent pattern and that of light transmitted through the field portion. In this regard, reference will be made to an example of sub peak resulting from interference in a mutually enhancing direction of light rays emitted from four positions, called quadrupole. An example of a pattern layout as seen in plan in quadrupole is shown in FIG. 20A and a section taken on line A–A′ is shown in FIG. 20B. In these figures, the reference numeral 100 denotes quartz glass (blank), numeral 101 denotes a halftone film, and numeral 103 denotes a circuit pattern. FIG. 14 shows an optical transfer image corresponding to the range from B to B′ in FIG. 20B. It is seen that the higher the transmittance of the halftone mask member, the larger the sub peak. To suppress this phenomenon, there is adopted a method wherein in the case of a high-transmittance halftone phase shift mask with a relatively high transmittance of a halftone film, a fine light shielding pattern is formed using Cr on a portion of the halftone film corresponding to a sub peak pattern generated position, thereby cutting off the transmitted light in that portion. This method is called a tri-tone halftone phase shift method because three tones are used which are a transmitting portion (glass portion), a halftone portion, and a light shielding portion (Cr portion).\nDescriptions on the tri-tone halftone phase shift are found, for example, in JP-A No.15130/1999, JP-A No.308715/1994, and JP-A No.90601/1997.\nFor providing a method of manufacturing an LSI of high integration and/or high-speed operation, it is necessary that a circuit pattern of fine dimensions be formed with a high dimensional accuracy. Therefore, a tri-tone halftone phase shift method having a high resolution is required to attain a high dimensional accuracy and a high phase control accuracy. In a halftone phase shift exposure method, exposure light which has passed through a halftone portion and exposure light which has passed through an aperture interfere with each other in the vicinity of a boundary between the aperture and the halftone portion to enhance the optical contrast and thereby improve the resolution and exposure latitude. Therefore, controlling the quantity of exposure light passing through the halftone portion, i.e., controlling the transmittance in the halftone portion, and phase control, are extremely important. Further, the pattern size accuracy of a halftone film greatly influences the dimensional accuracy of the pattern transferred. In a fine pattern near a resolution limit of a projection lens, the optical contrast decreases to a large extent due to optical diffraction, so that there arises a factor called MEF (Mask Error enhance Factor), which causes deterioration in dimensional accuracy of the transferred pattern to a greater extent than the dimensional accuracy of pattern on the mask. MEF is an index indicating to what degree a dimensional difference ΔLm of the transferred pattern is amplified relative to a dimensional difference ΔLw on the mask. It is represented by the following equation in which the reduction ratio of the projection lens is assumed to be M. For example, in case of using a 5× lens, the value of M becomes ⅕.MEF=ΔLm/(M·ΔLw)\nIn the case of such a fine pattern as used in a halftone phase shift mask, the value of MEF is usually in the range from 2 to 3, that is, variations in mask size are amplified from 2M to 3M and in this state there is made the transfer of pattern is made.\nRecent tri-tone halftone masks require a fine and extremely high accuracy, such as a pattern size on mask of 320 nm and a dimensional accuracy on mask of 9 nm. Therefore, the dimensional yield in this class is usually as low as 10% to 30%. That is, on the average, it is necessary that three to ten sheets of mask be fed for producing one sheet of a good mask. Thus, the cost of mask is high and the Turn Around Time (TAT) is low.\nAdditional factors such as the resist characteristic factor and the substrate step structure factor, the position and size of a fine light shielding pattern (auxiliary pattern) using Cr, which pattern is for preventing unnecessary image transfer caused by a sub peak, cannot be accurately anticipated by simulation. Therefore, even if a halftone phase shift film is satisfactory in point of dimensional accuracy, if it is found to be defective in the state of transfer after the fabrication of mask, it is necessary that the mask be fabricated again from the beginning. This has caused an increase in the TAT of the mask feed.\nIn the course of manufacturing a semiconductor device there are used about 30 sheets of photomask. One important factor of determining the semiconductor device development period is the supply of photomask, and for shortening the development period it is absolutely necessary to improve the photomask feed TAT. Further, the IC grade development is often conducted in wiring step and this grade development power also depends much on the feed TAT of the photomask for the wiring step."}
{"text": "This invention relates generally to personal computers, and particularly to a personal computer having a removable computer peripheral.\nPersonal computer systems in general and IBM personal computers in particular have attained widespread use for providing computer power to many segments of today's modern society. Personal computer systems can usually be defined as a desk top, floor standing, or portable microcomputer that consists of a system unit having a single system processor and associated volatile and non-volatile memory, a display monitor, a keyboard, one or more diskette drives, a fixed disk storage, and an optional printer. One of the distinguishing characteristics of these systems is the use of a motherboard, or system planar, to electrically connect these components together. These systems are designed primarily to give independent computing capability to a single user and are inexpensively priced for purchase by individuals or small businesses. Examples of such personal computer systems are IBM's PERSONAL COMPUTER XT and AT and IBM's PERSONAL SYSTEM/2 Models 25, 30, 35, 40, 50, 55, 56, 57, 60, 65, 70, 80, 90, and 95.\nThese systems can be classified into two general families. The first family, usually referred to as Family I Models, use a bus architecture exemplified by the IBM PERSONAL COMPUTER AT and other \"IBM compatible\" machines. The second family, referred to as Family II Models, use IBM's MICRO CHANNEL bus architecture exemplified by IBM's PERSONAL SYSTEM/2 Models 50 through 80. Many Family I models have used the popular INTEL 8088 or 8086 microprocessor as the system processor. These processors have the ability to address one megabyte of memory. Certain Family I and most Family II models typically use the high speed INTEL 80286, 80386, 80486, and Pentium microprocessors which can operate in a real mode to emulate the slower speed INTEL 8086 microprocessor or a protected mode which extends the addressing range from 1 megabyte to 4 Gigabytes for some models. In essence, the real mode feature of the 80286, 80386, and 80486 processors provide hardware compatibility with software written for the 8086 and 8088 microprocessors.\nPCs are particularly useful in a client-server environment where at least one PC system can be used as a server and other units, connected to the server, share the server's resources. Servers must often meet certain security requirements, possibly to notify the user of tampering. An unauthorized user may system may tamper with a system while running the authorized user's operating system (and application programs) and the authorized user may not be notified of the tampering until the next time the system is either powered off and then on or rebooted."}
{"text": "In recent years, significant efforts to suppress combustion noise and instabilities have been reported in the context of lean premixed gas turbine combustion systems. The effectiveness of an active combustion control system depends upon actuation, sensing, and control algorithms, among other factors. In spite of the significant progress in these areas, complete reliability of active combustion control in turbine engines is still a major concern, since an unexpected event can destroy the engine within a fraction of a second. Thus, passive techniques that do not affect the system volume or weight are desired to mitigate combustion instabilities and associated noise problems in turbine engines.\nIn view of the preceding, there is a need for a passive technique to suppress and/or mitigate combustion noise and/or stability."}
{"text": "This invention relates, in general, to amusement devices, and, in particular, to amusement devices with diverters that can be operated to change the direction of an opponent's ball."}
{"text": "Chrominance carrier modifiers are needed, for example, for reproduction of still pictures by PAL-color video signals recorded on magnetic tape in helical scan systems. In reproduction of still pictures, the tape advance is stopped, and the signal track is scanned by a magnetic head at 50 or 60 times a second. In a still picture reproduction of recorded PAL-color video signals, it is necessary to maintain the PAL-standard monotone alternation of the switching phase frequency which is lost when scanning a stopped tape.\nChrominance sub-carrier modifiers are known, see for example, the periodical \"Telefunken-Zeitschrift\", series 37, 1964, No. 2, pages 115-135. In such a known chrominance sub-carrier modifier, a reference carrier with double the frequency is modulated with a color signal F by means of a balanced modulator. The upper sideband is filtered out of the resulting heterodyne products. The remaining lower sideband F* corresponds to the conjugated complex form of the color signal F. An electronic switch thereafter selects out the proper sequence of the chrominance signal types F and F*.\nKnown chrominance sub-carrier modifiers with a double balanced or push-pull modulator have some disadvantages when used to reconstitute the phase relationships of PAL-video signals. Aside from the expense of the double balanced modulator and the frequency doubler, as well as the tolerance problems relating to the necessary phase adjustments, a shift from zero line offset to half line offset arises in the luminance component in the frequency range of the chrominance sub-carrier frequency if the original subcarrier is in a quarter line relationship. A reduction of horizontal resolution in one field is also observed. Disturbances due to color cros-talk arise as well.\nIt has previously been proposed--see British Pat. No. 1,015,179, Telefunken, to change signals following the standards of the PAL-color television system to those of a non-alternating system like the NTSC system but which is not adopted to the standardized NTSC system, in which special circuitry is provided to improve the color rendition even if phase errors may occur. As described, the color sub-carrier signal of one line, separated from the brightness signal, and the delayed color sub-carrier signal of the preceding line of substantially the same amplitude and phase--and also separated from the brightness signal--are added in a first combining stage and subtracted in another combining stage. The outputs from the two combining stages are added in an adder to the brightness signal from which the color sub-carrier signal has been removed. The negative polarity of the output signal in each alternate line is phase-shifted by 180.degree. by a phase shifter included in the circuit after the second combining stage, and operating at line frequency. The line-by-line alternation of the I-signal of the NTSC PAL-signal is thus eliminated. This signal--with reference to line alternation--is actually an NTSC signal."}
{"text": "1. Field of the Invention\nThis invention relates to a Beehive System that includes a small hive beetle trap and to a lure for capturing small hive beetles found in honeybee hives. It also relates to methods for using the system, the trap, and the lure.\n2. Description of the Related Art\nApiculture is a vital component of modern agriculture. Currently there are more then 2.5 million honeybee producing colonies in the United States with an estimated value of the beekeeping and pollination industries of about 14.6 billion dollars per year. Honeybees, specifically Apis mellifera L., are highly efficient foragers on numerous flowering plants, including agricultural crops. Plant pollination by the honeybee is the most commonly recognized means of crop pollination. Beekeepers are challenged with diseases and pests within hives such as the tracheal and varroa mites in the United States and destructive diseases such as foulbrood disease. A beehive generally has fixed top and bottom and includes an exclusion device to keep the queen bee isolated. The bottom board is generally a rectangular shaped box or frame upon which the rest of the beehive is stacked. The bottom board includes a gap on or in one of the four upstanding walls through which bees enter and exit the beehive while keeping out other larger insects, mice, and other rodents.\nThe small hive beetle is an exotic pest that originates from South Africa and was found in Florida in 1998. It has now spread throughout the eastern and mid-western United States, causing considerable damage to honeybee colonies and negatively impacting the beekeeping industry (Morse and Calderone, Bee Culture, Volume 128, 1-15, 2000). In sub-Saharan Africa, the beetle's native home, they are not considered a major pest of honey bees (Lundie, Sci. Bull. 220, Union of South Africa, Department of Agriculture and Forestry, 30 pages, 1940), because in strong colonies, worker honey bees encapsulate groups of small hive beetles with propolis (Neumann et al., Naturwissenschaften, Volume 88, 2140216, 200). However in the United States, where honey bee colonies are predominantly of European origin, small hive beetles have been reported to damage both weak and strong colonies (Sanford, Apis, Volume 16, 1-5; 1998). Defense mechanisms have yet to be demonstrated in European bees.\nDamage to honey bee colonies in the United States is caused mainly by the beetle larvae, which feed on honey, pollen, and brood. In addition, larval excrement tends to ferment honey, rendering it unsuitable for human consumption. In highly infested colonies, where larval feeding is extensive, bees generally leave. The beetle can multiply to huge numbers within infested colonies where it eats brood, destroys combs, and if uncontrolled ultimately destroys the hives. If the infestation is sufficiently heavy, the beetles may cause bees to abandon their hive. In stored combs and honey awaiting extraction, the beetle larvae may tunnel through the combs of honey, feeding and defecating, causing discoloration and fermentation of the honey. The small hive beetle will thrive in tropical, semi-tropical and temperate climates. It is an opportunist and a scavenger, taking advantage of situations that suit its reproduction. Adult beetles lay eggs in irregular masses in crevices or brood combs containing pollen or brood. Each female beetle is capable of laying an enormous number of eggs which can lead to a heavy infestation in a single generation. Beetle larvae mature from egg to pupae in about 10-16 days. The larvae require a diet of honey and pollen to fully develop. Beetle larvae also feed on stored comb and supers of honey in honey houses, especially frames containing pollen. The beetle is a major problem when honey supers are allowed to remain in the bee hive for long periods prior to extraction. Wax cappings, set aside following honey extraction, may also become beetle infested. Bee brood rearing stops when beetle infestation reaches high levels in the colony.\nCurrently small hive beetles are controlled using good management and a number of techniques that can reduce the impact of the beetle. Relocating hives may break the beetle's life cycle. Pesticide containing products such as CheckMite+™ strips and GardStar® (MannLake Ltd., Hackensack, Minn.) have been recommended for hive beetle control. The CheckMite+™ is an inside-hive treatment while the GardStar® is a soil drench.\nU.S. Pat. No. 6,468,129 (Griffith, Oct. 22, 2002) discloses a bottom board that has a screen permanently installed in the bottom board that covers the full expanse between the upstanding walls so that mites in the hive have to pass through the screen and conic in contact with a sub-board that is removable from the bottom board through a rear access. The subboard may optionally include a sticky coating to trap the mites that pass through the screen.\nU.S. Pat. No. 4,682,380 (Martin, Jul. 28, 1987) discloses a bottom board which is flat at its base with three upstanding walls to support the remainder of the beehive with a gap provided between two of the upstanding walls to allow bees to enter and exit the hive.\nSuazo et al., (Apidologie, Volume 34, 525-533, 2003) report that there has been no attempt to identify the cues attracting small hive beetles to honey bee hives. They state that Elzen et al. (Apidologie, Volume 31, 361-366, 1999) reported that small hive beetles are attracted to a combination of honey, pollen, and adult bees. They also reported that small numbers of beetles were captured in traps baited with adult bees but not in traps baited with honey and pollen, or brood alone.\nVarious forms of small hive beetle control are known, however there remains a need in the art for reliable control methods and traps for controlling small hive beetles. The present invention, as described below, is different from related art control methods, lures, and traps."}
{"text": "In wireless communication systems, there are two types of network topologies based on different paths between a source and a target. In the first type of network topology, the source and the target can communicate with each other directly, an example for the first type of network topology is the cellular network topology, which is as shown in FIG. 1. In the second type of network topology, the source and the target may communicate with each other by data forwarding of a Relay Node (RN). Networks with the second network topology are generally called forwarding networks below. Relay networks, Ad hoc networks, mesh networks, etc. related in the current communication field all belong to the forwarding networks.\nBased on the type of the data route, the forwarding network can be categorized as tree-like topology, mesh topology, hybrid topology, etc. The network with the hybrid topology is as shown in FIG. 2. According to the RNs' functions, RNs may be categorized as analog RN and digital RN. Furthermore, the digital RN includes active RN, half-active RN and passive RN.\nThe analog RN simply receives, amplifies signals, and then transmits the signals.\nThe active RN receives signals, processes the signals totally according to the control of the active RN, and then transmits the signals processed.\nThe half-active RN receives signals, processes the signals according to the control of the half-active RN and a controlling node of the half-active RN at the same time, and then transmits the signals processed.\nThe passive RN receives signals, processes the signals totally according to the control of a controlling node of the passive RN, and then transmits the signals processed.\nAs for the data transmission in forwarding networks, a typical manner is the hop-by-hop relay manner. That is, the data from super ordinate node is received via air interface and processed (the processing method may be different based on the different category of the node) by the current node, then the data is transmitted from one or more antennas of the current node and received by the subordinate node. As shown in FIG. 3, a forward network includes a Source Node (SN), three RNs (RN1, RN2, and RN3), and a plurality of Target Nodes (TNs). When the SN transmits data to the TN2, the data is transmitted through four hops (SN<->RN1, RN1<->RN2, RN2<->RN3, RN3<->TN2).\nIn forwarding networks with hop-by-hop relay manner, each RN occupies different radio resources (including time, frequency, codeword, space, and so on) for transmitting and receiving data. With this simple hop-by-hop relay manner, the system performance is improved, such as enlarging the coverage area, improving the user traffic, or increasing the system capacity, at the cost of occupying much more radio resources. Because each RN should occupy different radio resources for receiving and transmitting data, and neighboring nodes (or nodes within interference range) should use different radio resources to avoid degrading system performance incurred by increased system interference.\nThe prior art provides a forwarding node with a plurality of antennas. However, the inventor found in the inventing process that, the forwarding node in the prior art can not recognize the corresponding relationship between the antennas and the nodes in the coverage area of the antennas. Therefore, upon receiving data, the forwarding node in the prior art transmits data to designated areas by all antennas. Although the coverage area can be enlarged, user traffic can be improved, and the system capacity can be added when using this manner, a great amount of radio resources are still occupied. Moreover, a plurality of antennas transmit the same data will increase system interference in the forwarding network."}
{"text": "1. Field of the Invention\nThe present invention relates to a computer program product, system, and method for using views of subsets of nodes of a schema to generate data transformation jobs to transform input files in first data formats to output files in second data formats.\n2. Description of the Related Art\nAn Extensible Markup Language (XML) schema is a description of elements in a document encoded using the XML language. The XML schema may also be used to define a data model, also known as a meta model comprising a definition of a hierarchical representation of nodes that represent the data and relationship of data content for a data model. XML schemas that define more complex data models may have 1000s of nodes arranged in a complex hierarchy of trees of nodes. A large scale schema may be comprised of multiple XML Schema Definition (XSD) files and provides a definition of large scale trees of nodes and data content.\nFor instance, the financial industry uses FpML (Financial products Markup Language), which is an XML meta model providing a message standard for the financial industry, OTC Derivatives. The FpML meta model, version 4.9, groups more than 100 sub requests within 37 XSD files. Another large scale schema is ACORD family of XML specifications for Life, Annuity and Health that enables information exchanges among business partners, and one type in the schema has more than hundreds of sub-types.\nThere is a need in the art for improved tools for managing access of very large scale schemas and manipulate the data defined by such schemas."}
{"text": "The present invention relates to semiconductor MOSFET technology and more particularly to a trench MOSFET having self-aligned features.\nPower MOSFETs (metal oxide semiconductor field effect transistors) are well known in the semiconductor industry. One variety of power MOSFETs is the vertically-conducting trench MOSFET. A cross-section view of such a MOSFET is shown in FIG. 1. MOSFET 100 has trenches 111 each including a polysilicon gate 112 insulated from body regions 114 by a gate dielectric 110. Source regions 116 flank each side of trenches 111. Dielectric layer 120 insulates gates 112 from overlying metal layer 126. Substrate region 102 forms the drain of MOSFET 100.\nWhen MOSFET 100 is biased in the on state, current flows vertically between source regions 116 and substrate 102. The current capability of MOSFET 100 in the on state is a function of the drain to source resistance (Rdson). To improve the current capability of the MOSFET, it is necessary to reduce the Rdson. One way to reduce the Rdson of the trench MOSFET is to increase the trench density (i.e., to increase the number of trenches per unit area). This may be achieved by reducing the cell pitch. However, reducing the cell pitch of MOSFETs is limited by the particulars of the MOSFET cell structure and the specific process recipe used to manufacture the MOSFET. Reducing the cell pitch is made further difficult by such limitations of the manufacturing process technology as the minimum critical dimensions the photolithography tools are configured to resolve, the minimum required spacing between different cell regions as dictated by the design rules, and the misalignment tolerances.\nThe different dimensions that determine the minimum cell pitch for trench MOSFET 100 are shown in FIG. 1. Dimension A is the minimum trench width the photolithography tools are configured to resolve, dimension B is the minimum contact opening the photolithography tools are configured to resolve, dimension C is the minimum trench-to-contact spacing dictated by the design rules, and dimension D is the contact registration error tolerance or contact misalignment tolerance. The minimum cell pitch for MOSFET 100 thus equals A+B+2C+2D. Reduction of any of these dimensions without complicating the process technology is difficult to achieve.\nThus, a new approach wherein the cell pitch of the trench MOSFET can be reduced without increasing the process complexity is desirable."}
{"text": "This invention relates generally to systems for authenticating articles, methods for authenticating articles, and processes for marking articles for later authentication. This invention relates particularly to mechanisms and methods for detecting indicia which may include graphic images and/or characters and for comparing the detected indicia with predetermined indicia to verify the articles\"\" authenticity if the detected indicia are recognized. This invention relates more particularly to detection, discrimination, and recognition of fluorescent indicia rendered detectable by illumination of articles with ultraviolet light.\nThe counterfeiting of articles of many kinds has become a serious problem worldwide, causing great loss of revenues to legitimate businesses and to individuals. Counterfeiting has produced articles that are very difficult to distinguish from the genuine articles, taking sales revenues from the producers of genuine articles and impacting legitimate business\"\" reputations when the counterfeit articles have inferior quality and/or non-existent manufacturers\"\" support. Similarly, problems and losses occur due to counterfeiting of articles used for financial transactions or identification, such as credit cards, drivers\"\" licenses, passports, and immigration documents. Even when genuine articles have been marked with authentication markings such as watermarks, special papers, and holograms which are difficult to replicate, counterfeiters have managed to produce articles that appear genuine. The general availability of newer replication technology such as high-resolution image scanners, laser copiers and printers, and color-accurate color copiers and printers has enabled counterfeiters to produce more credible counterfeit articles. These technological developments of replication hardware continue to make rapid progress, as do the computer hardware and software that support them. Thus it is useful to have additional ways to mark articles for authentication with greater security. Greater security may be achieved by making authentication markings more difficult for counterfeiters to detect and interpret, by incorporating greater complexity into the markings, and by making replication by counterfeiters more difficult. Combining multiple kinds of marking indicia can further increase the complexity of detection, interpretation, and replication, thus providing even better security. Considering the security issue from another point of view, it is not desirable to use the same means that generated the visible appearance of an article to authenticate it, as the visible appearance is becoming easier to replicate. Invisible indicia that cannot be reproduced by common visible graphic copying and printing means can provide improved security, especially when combined with complex images, patterns, and/or characters.\nIn this specification, the following terms or abbreviations are used with the meanings listed here:\nCCDxe2x80x94charge coupled device\nDSPxe2x80x94digital signal processor\nexcitation timexe2x80x94a characteristic time for a fluorescent emission to be excited after ultraviolet illumination of a fluorescent substance begins\nextinction timexe2x80x94a characteristic time for a fluorescent emission to decay exponentially from its initial emission intensity IO to intensity IO/e\nIRxe2x80x94light in the infrared portion of the spectrum\nOCRxe2x80x94optical character recognition\nregistrationxe2x80x94the conventional definition of this term in the fields of graphic communication and printing, as defined, for example, in xe2x80x9cGlossary of Graphic Communicationsxe2x80x9d compiled by Pamela Groff and published by Graphic Arts Technical Foundation (Prentice-Hall, Inc., Upper Saddle River, N.J., 1998) at pages 245-246.\nUVxe2x80x94light in the ultraviolet portion of the spectrum\nvisible/IRxe2x80x94light either entirely in the visible portion or entirely in the infrared portion, or partly in both visible and infrared portions of the spectrum\nvisible/UVxe2x80x94light either entirely in the visible portion or entirely in the ultraviolet portion, or partly in both visible and ultraviolet portions of the spectrum\nMany methods have been known to authenticate valuable articles. Some known methods include imprinting a white-light hologram or imprinting reflective and diffractive indicia which display distinctive images that are difficult to counterfeit. Other known methods include incorporation of distinctive fibers into the articles, such fibers being detectable by visual observation, microwave irradiation, or other means. U.S. Pat. No. 4,921,280 describes fibers made luminescent by a dyeing process employing rare-earth compounds, which fibers may be incorporated into articles requiring authentication.\nIt is known that documents may be authenticated by marking the documents with substances such as inks or dyes that appear invisible or relatively unnoticeable to the naked eye in ordinary visible illumination, but that fluoresce when illuminated with ultraviolet light, revealing marks that serve to identify the legitimate document. These methods depend on substances that are not easily or inexpensively identified by a counterfeiter, and not easily or inexpensively duplicated or mimicked by a counterfeiter. When using these methods, it is desirable to use substances such as dyes or inks that fluoresce in narrow spectral bands which are distinguishable by sufficiently narrow-band detectors, but not readily distinguishable by eye. In U.S. Pat. No. 4,146,792 by Stenzel et al., these methods are extended to include dyes containing rare-earth elements whose fluorescence is influenced by the chemical environment of the fluorescing atoms in a non-fluorescing matrix, and the detection is refined to include detection of predetermined fine structure in the line spectrum of emitted light. Yet another class of authentication methods uses substances which fluoresce in the infrared portion of the electromagnetic spectrum when illuminated by light in the visible portion of the spectrum.\nIn U.S. Pat. No. 4,642,526 by Hopkins and assigned to the assignee of the present invention, a source of ultraviolet light is made self-modulating at a predetermined frequency. Detection of the secondary radiation, filtering of the detected signal, and demodulation of the filtered signal at the predetermined frequency allow the system of Hopkins\"\" invention to detect the fluorescent marks despite interference from ambient light sources.\nMarking products with indicia such as bar codes using fluorescent substances such as inks or dyes is also known in the prior art, both for the purposes described above and for providing identification on the products without detracting from the products\"\" appearance as normally viewed in visible light. U.S. Pat. No. 4,983,817 by Dolash et al. describes methods and apparatus for reading bar codes printed with fluorescent substances, while compensating for variations in background reflectance over the area printed with the bar code. U.S. Pat. No. 5,064,221 by Miehe et al. shows the use of fluorescent substances added to ink ribbon used for printing original documents, in order to distinguish the originals from copies. In many of the known authentication methods using fluorescence, the fluorescent identifying substance may be incorporated into the article during the article\"\"s manufacture. In U.S. Pat. No. 4,451,521 by Kaule et al., for example, luminophores are incorporated into paper during its manufacture. Some of the fluorescent substances used in the prior art for authentication purposes contain heavy metals, which are not as safe or environmentally benign as substances not containing such heavy metals. While fluorescent substances have been used in many ways in the prior art to mark articles with indicia for interpretation by human observers, to our knowledge neither optical character recognition (OCR) systems nor image verification systems have been combined with fluorescent markings and UV light illumination to provide improved authentication methods.\nThe technology of automatic pattern recognition is described in a large body of literature including numerous patents, especially in the application area of optical character recognition (OCR). Many of the recent developments in OCR technology (especially in recognition logic software) are usable with the present invention albeit with some adaptation to the type of light source and to such contrast conditions as may differ from those encountered with conventional OCR scanners having conventional visible/IR light sources.\nIt is an object of this invention to provide an improved authentication system which can detect and recognize complex authenticating indicia such as graphic images and alphanumeric characters. It is another object of the invention to provide an authentication system which can be readily automated to provide authentication that does not depend on human subjective judgment. It is another object of the invention to combine the use of substances which fluoresce under UV illumination, detectors capable of detecting fluorescent radiation, and image or character recognition techniques to enable the creation and authentication of articles with improved security. Another object of the invention is to enable the use of character or image indicia for authentication printed with substances that are safer and less likely to harm the environment than substances containing heavy metals. Another object of the invention is to enable the use of character or image indicia for authentication printed with substances that have sufficiently short extinction times to enable fast scanning. Another object is to enable capture of complex indicia without undue loss of resolution or smearing of details. Another object is to enable robust identification of articles of mail using mail bar codes such as xe2x80x9cPostNETxe2x80x9d codes printed with fluorescent substances. Another object of the invention is to enable the use of fluorescent indicia for authentication that are more complex than bar codes or simple arrays of dots and/or lines, and that allow a higher density of information per unit area than conventional bar codes. Another object is to provide an authentication system that makes counterfeiting and fraudulent labeling of articles including documents more difficult and expensive to accomplish. Another object is to provide a system that can mark articles for authentication without unduly compromising their aesthetic appearance. Another object is to provide authentication marks that are not easily removed. Another object is to provide flexible methods of authenticating articles that can be practiced in the field of use without requiring special printing press equipment. Another object is to provide an authentication system that can be re-programmed by a user to use various authentication criteria within its range of graphic image and character recognition capabilities.\nAn authentication system in accordance with this invention combines a source of ultraviolet light (and optionally a scanning mechanism) with apparatus for capturing and recognizing either graphic images or characters or both, where the graphic images and/or characters have been previously made with fluorescent substances that may be invisible under ordinary visible light, but are rendered detectable by the ultraviolet light. As in conventional optical character recognition (OCR), the characters may be conventional alphanumeric or other natural-language characters readable by human readers once they are made visible.\nThe authentication system has a housing enclosing its optical paths, a source of at least UV light, a detector for detecting graphic images or characters, conversion of the detector signal to digital form, a memory storing predetermined indicia, recognition logic, and indicators. The system may also include a scanning mechanism and optical filters to select predetermined wavelengths of fluorescent light. The detector may be capable of detecting both fluorescent images and normal visible images, and the authentication system may incorporate switching mechanisms to allow multiplexed acquisition of fluorescent and normally visible images produced by visible/IR or visible/UV illumination.\nThe authentication system can work in conjunction with processes performed in accordance with this invention, for marking articles with indicia selected from a predetermined set of graphic images and characters readable by optical character recognition and/or by image comparison. That process can include printing fluorescent graphic images or characters in registration (or in predetermined mis-registration) with substantially identically-formed images or characters previously printed with visible substances and/or with fluorescent substances of different characteristic fluorescent wavelength. The recognition logic of the authentication system can include comparison of these two (fluorescent and visible/IR, fluorescent and visible/UV, or two different fluorescent) images with each other. The process can also include printing fluorescent graphic images or characters with a multiplicity of fluorescent substances having distinct fluorescent wavelengths and excitation or extinction times, and printing those in overlay with registration or with predetermined registration offset. The registration offset may be selected to produce a predetermined stereo image when properly viewed by a human observer, thus allowing xe2x80x9cfirst orderxe2x80x9d authentication by visual inspection using suitable filters. Similar first order authentication may be provided by printing graphic image or character indicia which are both fluorescent and partially visible without UV illumination. Another overlay method uses two graphic images or character indicia, each incomplete in itself, with unconnected but complementary segments. When the two complementary segments are combined with appropriate registration, a complete recognizable graphic image, word, or natural language phrase is formed.\nOther objects, uses, advantages and features of the invention will become apparent to those skilled in the art from the following descriptions of preferred embodiments."}
{"text": "1. Field of the Invention\nThe present invention relates to an image forming apparatus, such as a copying machine, a printer, a facsimile machine, and a composite machine having functions of copying, facsimile, and the like.\n2. Description of the Related Art\nConventionally, as described in Japanese Unexamined Patent Publication No. HEI 04-356062, there exists an image forming apparatus comprising a drum unit including a photoconductive drum and a developing unit including a developing sleeve, the respective units being constructed as separate bodies and dismountably mounted to a main body.\nSuch image forming apparatus as described above comprises a chamber for dismountably mounting the respective units to a main body of the apparatus. On each of the opposite inner side walls of the chamber, guiding members to which guided members formed on the opposite exterior side portions of the respective units are engaged are formed so that the respective units are guidedly inserted to a predetermined mounting position in the main body.\nHowever, such conventional image forming apparatus has a problem that operability for insertion of the respective units is not so good since it is difficult to distinguish an initial position of an insertion path of each guiding member in the case where the guided portion of each unit is engaged with the respective guiding portion in the chamber for replacement and such. Especially, improvement in operability for insertion of the developing unit is requested since the developing unit is frequently mounted or dismounted with respect to the image forming apparatus for fixing up jamming of paper sheets.\nFurther, it is requested to prevent the drum unit and the developing unit from being inserted in wrong order."}
{"text": "This invention relates generally to vision protection devices, and more particularly, to a lens system for dentists, surgeons, and other practitioners which system is adapted to provide protection from ultraviolet radiation during surgical and diagnostic procedures.\nThe retina of the human eye is a thin structure of extreme complexity. The receptors, which may number 125,000,000 in each eye, comprise a large number (about 95%) of rods and a small number (about 5%) of cones. In a rod-free depression of the retina known as the fovea centralis, a small central bouquet of about 2,000 cones are packed together into a hexagonal array having a density of about 150,000 per square millimeter. The cones in this area are capable of dissecting the finest details of the optimal retinal image. The rods and cones are synaptically connected to the bipolar cells, which in turn relate to the ganglion cells, whose axons constitute the optic nerve fibers. There is a high degree of convergence, however, in that the 125,000,000 receptors ultimately feed into only 1,000,000 nerve fibers of the flexible optic nerve, which therefore constitutes the principal bottleneck of information flow in the visual system. The convergence ratio for the fovea is about 1:1, helping to preserve the high detail vision of this region, while in the peripheral region of the retina, the ratio is many thousands to one, leading to high sensitivity at the expense of resolving power.\nWhen the sensory layer of the retina separates from the pigment layer, retinal detachment is said to occur. This may result as a complication of some disease, or may be caused by repeated exposure to ultra violet radiation (light having a wavelength of between 10-380 nm). The detachment is partial at first, but without medical attention, almost always becomes complete, resulting in total and permanent blindness in the affected eye. At first, the patient may \"see\" flashes of light. Then, the patient may experience the sensation of a curtain gradually moving across the eye. The field of vision becomes progressively cloudy, until vision is lost. The progressive nature of retinal detachment is due to the gradual seepage of fluid from the large vitreous cavity into the space between the two layers of the retina. As more fluid seeps through the original hole in the retina, more of the sensory layer of the retina is separated from the pigment layer, until the detachment is complete.\nAs indicated above, repeated exposure to ultra-violet radiation can lead to retinal detachment and, if untreated, eventually to blindness. Accordingly, it is important to protect the eyes whenever and wherever exposure to sources of UV radiation is anticipated. Tanning booths, for example, typically provide their patrons with special eye pieces which fit over the eyes and filter out the ultra-violet radiation emitted therein. Because of the intensity of light used in such booths, the eyes of the patron would otherwise be exposed to UV radiation even when closed.\nLight sources which emit radiation in the UV range are also utilized in many dental and medical procedures. For example, UV light is typically employed to cure the epoxies or cements used in dental reconstruction or implantation procedures. During surgical or operating room procedures, in which localized, high intensity light is required, a \"cold light\" source is often employed. Essentially, the cold light source comprises a fiber optic cable which is illuminated by a high intensity high wattage lamp. The light source of the lamp is positioned to illuminate an inlet port of a fiber optic cable, which cable may be a few feet or more in length. The surgeon or dentist uses the fiber optic cable as a light source and he can therefore direct the light emanating from the outlet end of the fiber optic cable to any desired position. Because the amount of light needed for dental as well as surgical procedures is substantial, lamp sources such as halogen, or xenon or other high intensity lamps are employed. Each of these sources emit a sizable power in continuous ultraviolet radiation.\nFrom the foregoing discussion, it can be seen that it in many situations dental or medical practitioners are exposed to significant levels of UV radiation. For this reason, the practitioner must wear glasses or some other visual aid to filter out the harmful radiation. One disadvantage of such visual aids, however, is that they may diminish the clarity or otherwise obscure the practitioner's view of area of the patient under inspection. Accordingly, such aids are often removed when the practitioner desires a clearer view. A further disadvantage of such visual aids is that they do not accommodate the use of magnifying attachments which aid the practitioner in focusing on those areas of particular interest during a given surgical or diagnostic procedure.\nAccordingly, it is an object of the present invention to provide a visual aid device for use by dental or medical practitioners which device adequately protects the eyes of the practitioner from the harmful effects of ultraviolet radiation but which also permits the practitioner to obtain an unfiltered view when desired.\nIt is a further object of the present invention to provide a visual aid which is adapted to incorporate attachments having magnifying or other optical properties so that the device may be employed during typical surgical or diagnostic procedures."}
{"text": "The present invention is directed to a gas laser of the type having a vacuum-tight sealed housing which has a capillary extending at least partially in the housing and containing a discharge channel. Magnets are located in the proximity of the capillary, the magnetic field of the magnets penetrating the discharge channel in the capillary at least over a portion of its length. Such a prior art gas laser is disclosed in U.S. Pat. No. 4,035,741. In the prior art gas laser, magnets are located outside a housing which is sealed so that the housing is vacuum-tight and are located in the proximity of a capillary that projects over the housing. A relatively complex adjustment mechanism is required in order to adjust the magnets relative to the capillary. The arrangement disclosed basically covers only one part of the capillary.\nIBM Technical Disclosure Bulletin, Vol. 22, No. 8A, January, 1980 discloses a similar arrangement in which a row of magnets is located in proximity to the discharge tube."}
{"text": "(i) Technical Field\nThe present invention relates to methods for fabricating semiconductor devices, and more particularly, to a method for fabricating a semiconductor device having a silicon nitride film formed on a nitride semiconductor layer.\n(ii) Related Art\nSemiconductor devices using nitride semiconductors are used in, for example, power devices capable of operating at high frequencies and outputting high power. Particularly, FETs (Field Effect Transistors) such as HEMT (High Electron Mobility Transistor) are known as transistors suitable for amplifying in high-frequency or RF bands, which may include microwaves, quasi-millimeter waves, or millimeter waves.\nIt is known to form a silicon-rich silicon nitride film formed on a nitride semiconductor layer in order to suppress the collapse phenomenon of the drain current (see Japanese Patent Application Publication No. 2006-278812.\nHowever, a high composition ratio of silicon in the silicon-rich silicon nitride film results in a drift of the drain current."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for converting a standard text file into a text file for a xe2x80x9cfill-in-the-blanksxe2x80x9d type of teaching method and apparatus which can be implemented on a computer for learning information using the converted text file.\n2. Description of the Related Art\nComputer technology has been rapidly advancing, and personal computers have recently been used as teaching and learning tools. There are already a number of prior art systems which display xe2x80x9cfill-in-the-blanksxe2x80x9d questions. However, no prior art method has been directed towards how to easily create a data file for a xe2x80x9cfill-in-the-blanksxe2x80x9d type of question and answer program.\nU.S. Pat. No. 5,885,083 issued Mar. 23, 1999 to Ferrell, and discloses a method and system for interactive speech training. This invention includes Sentence Completion and Story Completion. According to the disclosed method, sentences with an omitted word or words must be stored in memory. Further, the correct answer corresponding to the omitted word or words and incorrect answers must also be stored.\nU.S. Pat. No. 5,810,599 issued Sep. 22, 1998 to Bishop, and discloses a system and method for interactive audio-visual foreign language maintenance. This invention utilizes a xe2x80x9cfill-in-the-blanksxe2x80x9d method for teaching foreign languages to students. According to this method, it is necessary for the system to store two kinds of sentences. One sentence comprises lines of dialog with selected word left blank, and the other sentence comprises the completed lines of dialog.\nIn both of these references, two kinds of data must be stored in a file or files on a storage medium. Accordingly, a user must edit both kinds of the data when creating or updating them. This is so troublesome that a dedicated program for editing the data is necessary.\nFurther, both prior art references require a user to input an answer. If a user inputs an answer, the computer can evaluate the answer and indicate whether the user input is correct. However it is time-consuming for a user to input an answer using a keyboard or other conventional input devices, especially when the questions are numerous.\nHTML files and browsers such as Netscape and Microsoft Internet Explorer are well known in the art. A xe2x80x9clinkxe2x80x9d which is used in a HTML file is also well known in the art. A link to another place in the same document looks like this:\n less than A HREF=xe2x80x9c#targetxe2x80x9d greater than Click here. less than A greater than \nIf the user clicks on the words xe2x80x98Click here.xe2x80x99, the browser will jump to the location in the current document marked with  less than A NAME=xe2x80x9ctargetxe2x80x9d greater than  This is the target of above link. less than /A greater than \nUsing a link, a data file for a xe2x80x9cfill-in-the-blanksxe2x80x9d program can be made as follows:\nThe  less than 1 greater than  of the invention, which should be as short and specific as possible, should appear as a  less than 2 greater than  on the  less than 3 greater than  page of the specification, if it does not otherwise appear at the beginning of the application.\n1=title\n2=heading\n3=first\nHowever, a user will have difficulty in converting an ordinary text file into the data file like this. A user must extract the keywords, store them in an appropriate location, and link them. In order to do this, therefore, a dedicated program is necessary.\nNot only text but also images can also be displayed. It is known that images often gives a student a strong impression, and therefore enhances the learning experience. The prior art discussed above displays images. However the position where images are displayed is fixed. After the above invention\"\"s systems are completed, the display layout is not movable.\nBelow is a example of an HTML file.\n less than HTML greater than \n less than HEAD greater than \ntriangle\n less than /HEAD greater than \n less than BODY greater than \nThis is a regular triangle. less than BR greater than \n less than IMG SRC=xe2x80x9ctriangle.jpgxe2x80x9d WIDTH=xe2x80x9c60xe2x80x9d HEIGHT=xe2x80x9c40xe2x80x9d greater than \n less than /BODY greater than \n less than /HTML greater than \nUsing an HTML file, both text and images can be displayed in a window on the display device. The positions where images are displayed are not fixed. The position where images are displayed can be altered by altering the HTML file. However it is not visually clear for a user where and how large images are displayed when a user views the HTML file.\nU.S. Pat. No. 5,141,439, issued Aug. 25, 1992 to Cousins, discloses a keyword teaching and testing method. According to this reference, a student does not use a computer but uses printed matter, something to write with, and something to delete keywords with, such as typewriter correction fluid. It is difficult and time-consuming to delete keywords physically especially when printed characters are small.\nIt is therefore an object of the present invention to provide a method and system for learning information by keyword methodology using a computer which enables a user to easily convert a text file into a text file for a xe2x80x9cfill-in-the-blanksxe2x80x9d question program.\nAnother object of the present invention is to provide a method and system for presenting already acquired information through as easy operation as possible.\nA still further object of the present invention is to provide a method and system for presenting images with text in a window on a display device which enables a user to visually control and understand the position and size of the images when editing the text file.\nThe invention, therefore, includes a method of creating a computer-based teaching system, with the method comprising modifying a text file by inserting a first predetermined character at a first end of a first selected information element, and inserting a second predetermined character at a second end of the first selected information element. A blank box is created in an area defined by the first and second predetermined character. The blank box replaces the first selected information element. Data is stored in memory, representing the first selected information element. A user interface is provided, which enables selective display of either the blank boxes or the first selected information element.\nThe invention also includes an apparatus for teaching information. The apparatus comprises a display device, and an input device for inputting commands of a user. A computer is connected to the display device and the input device, with the system including an editing unit for editing a text file, and a storage unit for storing the text file. A position determining unit is provided for determining selected positions of predetermined characters with respect to selected information elements in the text file. A converting unit converts the text file into a combination of text and selected blank boxes based upon the selected information elements."}
{"text": "A semiconductor memory device in semiconductor devices may store digital data. As the electronic and semiconductor industries become highly sophisticated, the requirement for high-integration of semiconductor memory devices is gradually increasing. For example, as portable electronic devices such as laptop computers, mobile phones, digital cameras and MP3 players advance, requirements for semiconductor memory devices capable of storing more data are increasing. To satisfy these consumer requirements, highly integrated semiconductor memory devices are required.\nGenerally, the minimum line width of fine patterns constituting a semiconductor memory device may be decreased for higher integration of a semiconductor memory device. By two-dimensionally decreasing the minimum line width of the fine pattern, more memory cells may be integrated in a limited area.\nHowever, due to various factors such as photolithography process parameters, the amount by which the minimum line width is decreased may be limited. Moreover, as the line widths of the fine patterns decrease, the characteristics of the fine patterns may be deteriorated, thereby reducing semiconductor memory device reliability. Accordingly, much research is actively being conducted in pursuit of new methods for implementing semiconductor memory devices that are highly integrated and have superior characteristics."}
{"text": "The present invention relates to an improved method and apparatus for interconnecting large scale integrated (LSI) circuits and, more particularly, associative techniques for interconnecting LSI elements on a wafer to enable wafer scale integration.\nPresent methods of semiconductor fabrication typically require that a plurality of identical semiconductor elements or circuits be fabricated on a single wafer substrate by a series of process steps. At the conclusion of this fabrication process, the elements or circuits are tested and the defective elements identified. The wafer is then scored and diced into individual parts, each containing a complete circuit. Finally, the operable parts are packaged to provide external electrical connections and appropriate environmental protection.\nIt has long been realized that circuit costs as well as space and power requirements could be reduced if the operable parts or elements on the wafer could be interconnected on the wafer itself. This would reduce the necessity for circuit boards and interconnection wiring between elements, which would, in turn, result in a decrease in space and cost.\nA variety of methods for interconnecting elements on a wafer have been suggested. One method is to add another metalization step after the circuits or elements of the wafer are tested to interconnect the operable element on the wafer. The difficulty with this method is that a different mask is required for each individual wafer, greatly increasing the cost of the finished product.\nAnother method is the inverse of the discretionary wiring concept attempted in the 1960's to achieve LSI circuits with small scale integrated and medium scale integrated circuits. Usually, the arrays to be connected are storage arrays. Each array consists of a storage portion, an address mechanism for accessing the data and a permanent disconnection mechanism. The permanent disconnection mechanism is accomplished by blowing a fuse, or charging up a floating gate MOS (FAMOS) device, or laser burnout, etc. The disconnection isolates a defective array from the wafer bus so that the defective arrays do not interact adversely with the operative arrays. Although this method is suitable for wafer scale integration, it is limited in application. It requires a hardwired decoder and does not permit reconnection of the spring of arrays for increased reliability.\nAnother wafer integration scheme, as described in U.S. Pat. No. 3,940,740, provides for spare rows and columns on the wafer matrix for appropriate sparing in case of defective elements. One disadvantage of this configuration is that an entire row of elements must be provided, in the worst case, to spare a single defective matrix element. A preferred method would be one where spare elements can be inserted into the matrix on a random basis to allow all spare elements to be used.\nIn addition, all of these methods require extensive decoding circuits and uniquely dedicated decoding bus lines for element enable, which reduce the usable area of the wafer and add possible failure sites. Because of these reasons, wafer scale integration is seldom achieved.\nWhat is required is a wafer scale integration configuration wherein a relatively small number of spare elements and a relatively simple decoding circuit are employed to obtain high reliability and high functional circuit density at the wafer level.\nIt is also an object of this invention to allow the replacing circuit element to assume the address of the replaced element so that the system need not be reprogrammed after the replacement."}
{"text": "Systems for sequentially feeding individual open bags of a granular product such as seed or animal feed to the head of a sewing machine after placing a tag at the mouth or opening of the bag to both close the mouth of the bag and attach the tag to the mouth of the bag through a single stitching action are known in the packaging industry. Examples of two such systems are disclosed in U.S. Pat. No. 3,805,716 of Cerioni and U.S. Pat. No. 5,479,757 of Ogawa. These two prior art systems are each arranged to deliver a single tag or label for each bag being conveyed to the sewing machine. The Fischbein 3920 dual tag placer is capable of feeding or conveying two tags or labels onto each bag for simultaneous connection thereto in the same bag-closing stitching action. However, the tags are placed side by side on the bag and therefore may occupy a significant fraction of the surface area of the face or side of the bag on which they are placed. It some cases, it is desirable to leave an area open for receipt of other labeling or marking directly on the bag, for example by ink printing onto the bag itself.\nApplicant is unaware of any prior developments in this field facilitating automated feeding of more than two tags for attachment to a single bagged product or facilitating feeding of two tags in a more space-efficient configuration than the machine identified above."}
{"text": "When color development processing is carried out after a silver halide photographic light-sensitive material is exposed to light, a developing agent, such as an aromatic primary amine, oxidized with silver halide reacts with a coupler to form a color image. In this process, color reproduction by a subtractive process is generally utilized. In accordance with this process, color images of yellow, magenta and cyan, which are complement colors of blue, green and red, respectively, are formed in order to reproduce images of blue, green and red.\nPhenol derivatives or naphthol derivatives are mainly used as cyan color image forming couplers. However, the color images formed from conventionally employed phenol derivatives or naphthol derivatives have some problems with respect to preservability. For example, color images formed from the 2-acylaminophenol cyan couplers as described in U.S. Pat. Nos. 2,367,531, 2,369,929, 2,423,730 and 2,801,171 generally have inferior fastness to heat. Color images formed from the 2,5-diacylaminophenol cyan couplers as described in U.S. Pat. Nos. 2,772,162 and 2,895,826 generally have inferior fastness to light, and color images formed from 1-hydroxy-2-naphthamide cyan couplers generally have inferior fastness to both light and heat (particularly humidity and heat).\nFurther, 2-acylaminophenol cyan couplers and 1-hydroxy-2-naphthamide cyan couplers are defective in that the color density obtained therefrom severely decreases when processed with a bleaching solution having weak oxidizing power.\nCyan couplers derived from m-ethylphenols having an .alpha.-phenoxyalkyl group as a ballast group as described in U.S. Pat. No. 2,772,002 exhibit remarkably improved heat fastness of color image and extremely small reduction in color density when processed with a bleaching solution having weak oxidizing power. However, these cyan couplers have a low color forming property and have a side effect of degradating light fastness of the color image.\nOn the other hand, cyan couplers derived from methylphenols having an alkyl group as a ballast group as described in U.S. Pat. No. 4,686,177 have a markedly high color forming property and are excellent in light fastness of the color image formed therefrom. However, the heat fastness of the color image and the decrease in color density when processed with a bleaching solution having a weak oxidizing power are inferior, as compared with the cyan couplers as described in U.S. Pat. No. 2,772,001 mentioned above."}
{"text": "Data transaction processing systems including multiple transaction processors typically have a finite or fixed amount of computing resources to be allocated amongst the various transaction processors. The computing resources may be of different types, e.g., input bandwidth compared to output bandwidth for a given transaction processor, or overall bandwidth, or time division frequency access, quality of service, or priority, associated with one transaction processor as compared to another transaction processor. Computing resources that can be allocated also include memory, CPU processing cycles, cores or threads dedicated to performing specific tasks. An efficient computing system allocates an optimum amount of its finite resources to the various processors, and the tasks performed by its processors. In a data transaction processing system including multiple different data transaction processors having finite computing resources, the optimum allocation of the various computing resources may vary for each of the data transaction processors."}
{"text": "Serotonin Selective Reuptake Inhibitors (SSRIs) currently provide efficacy in the treatment of several CNS disorders, including depression and panic disorder. SSRIs are generally perceived by psychiatrists and primary care physicians as effective, well-tolerated and easily administered. However, they are associated with a number of undesirable features.\nThus, there is still a strong need for compounds with an optimised pharmacological profile as regards the activity on reuptake of the monoamine neurotransmitters serotonin, dopamine and noradrenaline, such as the ratio of the serotonin reuptake versus the noradrenaline and dopamine reuptake activity.\nU.S. Pat. No. 3,196,154 (Sterling Drug Inc) describes 3-substituted-9-methyl-3,9-diaza-bicyclo[3.3.1]nonanes. The compounds 9-methyl-3-phenyl-3,9-diaza-bicyclo[3.3.1]nonane and 3-(3-chloro-4-methylphenyl)-9-methyl-3,9-diaza-bicyclo[3.3.1]nonane are disclosed.\nFales H M and Barnes R A [J. Am. Chem. Soc.; 1954; 76(7); 1947-48] describe the synthesis of 9-methyl-3,9-diazatricyclo[3.3.1.23,9]undecane. The synthesis of the compound N-picryl picrate of 9-methyl-3,9-diaza-bicyclo[3.3.1]nonane is described."}
{"text": "In order to balance a family's increasingly limited free time with necessary chores and responsibilities many people need frequent and specialized outside help in their households. Typically, a household might require assistance with duties, such as, cooking, cleaning, pet care, babysitting, shopping, errands, eldercare, lawn and garden maintenance or handyman duties. For many households, however, the need for such services is often not large enough to require regular full-time help. Thus, when such a need arises, members of the household are faced with the difficult problem of both locating and scheduling competent individuals to perform the required service.\nThe problems faced by households requiring part-time work are mirrored by personnel seeking to provide these services. The irregular needs of household clientele make it difficult for a service provider to connect with individuals requiring these part-time and specific services. Moreover, there exists a group of individuals, such as people with regular full or part-time jobs, that would be willing to perform these services if they could fit them within their schedule, but for whom the logistical issues raised in establishing clients make performing this type of extra work impractical."}
{"text": "This invention relates to a stroke enlarging mechanism of an arm and the like, and more particularly to a type thereof adapted to be used in a work removing device for a die-casting machine or a press.\nIn a technical field of, for instance, industrial robots and manipulators, it is frequently required that a chuck-hand thereof grasps and transports products or works from one position to another along a linear path, and such functions have been accomplished by an arm coupled to a piston operated by a hydraulic or pneumatic cylinder.\nIn the above described case, although there is no problem when the distance of the required movement of the arm is equal to or less than the stroke of the piston, a mechanism for enlarging or multiplying the stroke of the piston is required when the distance of the required movement is larger than the stroke of the piston.\nA typical construction of a conventional stroke enlarging mechanism is shown in FIG. 1. In this mechanism, a pinion 14 is mounted to be freely rotatable at an end of a piston rod 12 movable in a cylinder 10, and a pair of movable racks 16 and 18 are disposed so that the pinion 14 engages with the racks 16 and 18 at an upper end and a lower end thereof, respectively. When the piston-rod 12 moves outwardly from the cylinder 10, the rack 16 is driven upward while the rack 18 is driven downward as viewed in FIG. 1. Although an enlargement of the stroke of the piston can be realized by the above described stroke enlarging mechanism, this mechanism requires a space substantially equal to the enlarged stroke, and therefore it has been difficult to use the mechanism in a narrow space such as a factory of a low ceiling height or where the vertical space therein is limited because of the presence of a hoist or a crane, where the mechanism is disposed vertically, and in a narrow space horizontally where the mechanism is disposed horizontally."}
{"text": "1. Field\nThe disclosed embodiments relate to techniques for sharing data among users. More specifically, the disclosed embodiments relate to techniques for sharing and synchronizing data across users of cloud computing systems.\n2. Related Art\nNetwork-enabled users commonly use cloud computing systems to store, manage, and/or synchronize data. For example, a user may store images, audio, video, documents, and/or other files on an electronic device such as a laptop computer, personal computer, tablet computer, mobile phone, and/or portable media player. The user may also have a user account with a cloud computing system that allows copies of the files to be stored and/or backed up at the cloud computing system. The cloud computing system may additionally transmit the copies to other electronic devices of the user, thus synchronizing the data across the users' devices and with the user account.\nHowever, automatic synchronization and/or sharing of data are typically unavailable between users and/or user accounts of cloud computing systems. Instead, a user may share data with another user using conventional digital communication mechanisms, which may require manual setup and/or configuration by the user. For example, the user may share a photo album with the other user by uploading the photo album to the user's user account with a photo-sharing site and sending a link to the photo album to the other user via email, instant message, and/or text message. The other user may then access the photo album by directing a web browser to the link. As a result, the user may be required to select the appropriate link to the photo album and address the email, instant message, and/or text message to the other user to successfully share the photo album with the other user. Alternatively, the user may upload the photo album to a remote directory with the cloud computing system and manually grant access to the remote directory for the other user. However, sharing of the photo album through the cloud computing system may not work if the user uploads the photo to the wrong directory, shares the wrong directory with the other user, and/or shares the correct directory with the wrong user. Consequently, existing techniques for sharing and/or synchronizing data between users may be tedious, confusing, error-prone, and/or time-consuming for the users."}
{"text": "Prior Art\nM. B. Frankel, J. Org. Chem., 26, 4709 (1961).\nReference shows the reaction of 2,2-dinitro-1,3-propanediol with polynitroaliphatic primary amines and formaldehyde to give 1,3-bis(3',3',3'-trinitropropyl) and 1,3-bis(3,3-dinitrobutyl)-5,5-dinitrohexahydropyrimidine. PA1 Reference shows thorough investigation of the products of the reaction described in the previous reference as well as the nitration of 1,3-bis(3',3',3'-trinitropropyl)-5,5-dinitrohexahydropyrimidine to give 1,1,3,6,6-hexanitro-3-azaheptane. PA1 Reference describes the nitration of 7-nitro-1,3,5-triazaadamantane. Seven 1,3,5-trinitrohexahydropyrimidines, each with a different additional substituent at the 5-position were obtained.\nR. G. Gafurov, S. I. Sviridov and L. T. Eremenko, Izv. Akad. Nauk. SSSR, Ser. Khim., 6, 1366 (1973).\nA. F. Farminer and G. A. Webb, J.C.S. Perkin I, 940 (1976).\nMention may also be made of two compounds which find widespread usage as rocket fuel oxidants. The first of these compounds, commonly termed HMX, is 1,3,5,7-tetranitro-1,3,5,7-octahydrotetrazocine. This compound has the structure ##STR2##\nThe other compound, commonly termed RDX, is 1,3,5-trinitro-1,3,5-hexahydrotriazine. This compound has the structure ##STR3##\nAs compared to the DNNC compound of the present invention, the compounds of the prior art are relatively lean in oxygen in the sense of containing little if any oxygen over and above that required to burn the compound to CO and H.sub.2 O. In expressing the available oxygen value of a given oxidant compound, reference can be made to its \"oxygen balance\". This characteristic is determined by first calculating the number of oxygen atoms theoretically left in the compound after burning the same to CO, H.sub.2 O and N.sub.2 in the absence of any external oxygen source. The number so obtained is then divided by the molecular weight of the oxygen compound under test, with the value so obtained being multiplied by 1,000. Using this method all the prior art compounds mentioned above are found to have an oxygen balance equal to, or less than 0%, while DNNC has the relatively high oxygen balance of 6%."}
{"text": "Devices are known which use pulse-echo ultrasound to provide information regarding the thickness, configuration, composition and other characteristics of internal features of various parts of a body. Such a device is described, for instance, in U.S. Pat. No. 4,576,177 to Webster. However, specific problems are encountered when utilizing pulse-echo ultrasound to acquire information with sufficiently high resolution in extremely close range applications, such as acquiring information relating to the walls of coronary arteries.\nMore specifically, in known long range ranging systems, such as pulse-echo radar or sonar, a single electro-acoustical device is typically used for transmitting the sounding pulses and receiving the echoes. The device is switched between transmit and receive modes using a TR switch. The latter is designed to operate so as to minimize any interference during the receive mode caused by the excitation of the electro-acoustic device during the transmit mode. For example, the instant the device stops transmitting it may continue to ring. Accordingly, in order to receive any echo signals it is important that the device be allowed to stop ringing before the device is switched to the receiving mode. The TR switch can be operated quickly to quench the device at the termination of the transmit mode in order to switch the device to the receive mode. However, a finite time is nevertheless required to expire before the electro-acoustic device is sufficiently quiescent in order to operate in the receive mode without residual interference from the transmit mode. This finite time, called \"dead space\", thus creates a minimum time the receiver must be turned off following the end of the transmit mode. In other words, the target distance must be sufficiently large so that echoes are not received by the electro-acoustic device until after the time of the dead space has elapsed. The dead space therefore creates a minimum target distance at which the ranging device will operate without cross interference between the two modes and without a loss of information. This is usually not a problem for typical radar and sonar applications where target distances are well beyond the minimum distance required.\nIn an effort to improve short range performance, some short range pulse-echo ranging systems have resorted to using separate transmitter and receiver devices. The amount of excitation from the transmit pulse resident in the receive channel is naturally reduced by the physical isolation of the two (often by physically spacing one a relatively large distance from the other) and, hence, the switching need not be so complete as is the case when a single device is used to transmit and receive. Additional improvements can be achieved by utilizing improved materials which exhibit reduced ringing. For example, as described hereinafter, certain plastic materials have been found particularly useful in medical imaging applications because of the reduced acoustic impedance of such materials relative to known crystalline materials. However, such systems are still range limited by the required dead space, and the physical separation of the two transmitter and receiver devices. Therefore, until the present invention the use of separate transmitter and receiver devices for pulse-echo ultrasonic acquisition of information regarding very close range targets, such as the walls of coronary arteries, has been impractical.\nMore specifically, the system disclosed in the Parent Application comprises a catheter probe that is adapted to be inserted into a part of a body, and is particularly good at providing relatively high resolution imaging data of an relatively small, predetermined portion of the body, such as a small section of a coronary artery (the catheter may also be adapted to deliver laser energy to the interior of the body part for modifying internal features thereof, e.g., removing plaque deposits from a coronary artery). A transducer assembly is attached to the distal end of the catheter for emitting and receiving acoustic pulses used in generating imaging information. Sets of imaging data are created by moving the catheter axially along and notationally about its axis, within the body section of interest, through a series of imaging locations, while the transducer assembly is actuated to emit a train of acoustic pulses and responsively receive a series of acoustic echoes at each imaging location. By also continuously sensing the location of the catheter, and relating the set of data to the respective location where it was generated, an image of the internal features of the body part may be generated.\nIt is clear that the target being imaged by the system disclosed in the Parent Application is extremely close to the transducer assembly. In fact the transducer assembly is frequently almost in contact with the targeted surfaces, such as the inner surface of the wall of an artery. At such close range, relying on standard radar and sonar techniques of allowing the required time of dead space to transpire after the transmit mode and before switching to a receiving mode is inadequate since the target ranges are well below the minimum range for the dead space required and imaging data will be lost during this time. Utilizing two transducer devices, one for transmitting and one for receiving, does not alone overcome the problem since the two devices must be mounted on the end of a catheter extremely close to one another. Because of the close proximity between the two devices, sufficient dead space still must be provided to prevent acoustic interference, such as acoustic cross talk, between the transmitting device and receiving device, as for example when the transmitting device is ringing following the generation of an acoustic pulse. In fact, because of the close proximity of the two, certain electrical interference problems can occur between the transmitting and receiving devices which do not occur with a single device, such as capacitive and inductive coupling between the two devices."}
{"text": "1. Field of the Invention\nThis invention relates to catalyst system for the polymerization of olefins, particularly, to a catalyst system comprising a conventional supported Ziegler-Natta catalyst and a lithium cyclopentadienide or lithium indene compound used to polymerize .alpha.-olefins, such as propylene.\n2. Description of the Prior Art\nPolyolefin manufacturing processes typically involve the polymerization of olefin monomer with an organometallic catalyst of the Ziegler-Natta type. Catalyst systems for the polymerization of olefins are well known in the art. Typically, these systems include a Ziegler-Natta type polymerization catalyst component; a co-catalyst, usually an organoaluminum compound; and an external electron donor compound or selectivity control agent, usually an organosilicon compound. Examples of such catalyst systems are shown in the following U.S. Pat. Nos.: 4,107,413; 4,294,721; 4,439,540; 4,115,319; 4,220,554; 4,460,701; and 4,562,173; the disclosures of these patents are hereby incorporated by reference.\nA Ziegler-Natta type polymerization catalyst is basically a complex derived from a halide of a transition metal, for example, titanium, chromium or vanadium with a metal hydride and/or a metal alkyl, typically an organoaluminum compound, as a co-catalyst. The catalyst is usually comprised of a titanium halide supported on a magnesium compound complexed with an alkylaluminum co-catalyst.\nLithium cyclopentadienide and lithium indene compounds are known in the synthesis of metallocene compounds useful as catalysts. U.S. Pat. No. 5,168,111; 5,096,867; 5,057,475 and 5,055,438 disclosed using a cyclopentadienyl lithium compound in a synthesis of a bridged metallocene-alumoxane catalyst system for the production of polyolefins and ethylene-.alpha.-olefin copolymers. U.S. Pat. No. 5,120,867 disclosed use of lithium indene in the synthesis of a silicon-bridged metallocene.\nLithium compounds have been used in hydrogenation of conjugated diolefin polymers. U.S. Pat. No. 4,673,714 disclosed the use of hydrocarbon lithium compounds, specifically alkyl lithium, with bis(cyclopentadienyl)titanium compound to hydrogenate the unsaturated double bonds of conjugated diolefins and an improved process which eliminated the need for these alkyl lithium compounds.\nMolecular weight can be increased in polymerization processes using conventional supported Ziegler-Natta catalysts in several known ways.\nIt would be advantageous to increase molecular weight in polymerization processes using conventional supported Ziegler-Natta catalyst by the addition of a compound which interacts or modifies the catalyst without changing process condition or other catalyst components."}
{"text": "The present invention relates in general to surgical devices and procedures, and more particularly to minimally invasive surgery.\nSurgical procedures are often used to treat and cure a wide range of diseases, conditions, and injuries. Surgery often requires access to internal tissue through open surgical procedures or endoscopic surgical procedures. The term “endoscopic” refers to all types of minimally invasive surgical procedures including laparoscopic, arthroscopic, natural orifice intraluminal, and natural orifice transluminal procedures. Endoscopic surgery has numerous advantages compared to traditional open surgical procedures, including reduced trauma, faster recovery, reduced risk of infection, and reduced scarring. Endoscopic surgery is often performed with an insufflatory fluid present within the body cavity, such as carbon dioxide or saline, to provide adequate space to perform the intended surgical procedures. The insufflated cavity is generally under pressure and is sometimes referred to as being in a state of pneumoperitoneum. Surgical access devices are often used to facilitate surgical manipulation of internal tissue while maintaining pneumoperitoneum. For example, trocars are often used to provide a port through which endoscopic surgical instruments are passed. Trocars generally have an instrument seal, which prevents the insufflatory fluid from escaping while an instrument is positioned in the trocar.\nWhile surgical access devices are known, no one has previously made or used the surgical devices and methods in accordance with the present invention."}
{"text": "Recently, the demand for foldable futons has increased due to the flexibility of use in that these futons may function as either a sofa or a bed. The dual functionality of the foldable futons makes them both cost effective and space saving. Foldable futons abrogate the need to have a separate sofa and bed. To further decrease the cost of foldable futons, many are manufactured with easily workable and inexpensive materials. In this fashion, futons may be manufactured without the need for highly skilled labor. This further reduces the cost of futons, because they may be manufactured in under-developed countries in which the workers earn a lower annual wage than workers in more developed countries.\nCurrently, foldable futon frames are available in a variety of designs. For example, U.S. Pat. No. 5,303,432 to Fitts discloses a convertible futon frame that includes opposed end frames interconnected by front and back rails to define a rigid rectangular frame. Seat and back members are movably connected to the rectangular frame and to each other so that an individual can convert the frame between sitting and sleeping positions.\nTypically, foldable futon frames are made from either metal or wood. Metal futon frames are typically more durable than wooden futon frames. On the other hand, wooden futon frames are often more appealing, aesthetically, than metal futon frames. However, wooden futon frames suffer another drawback in addition to reduced durability, when compared to metal frames.\nAs discussed above, foldable futon frames are often produced in under-developed countries resulting in the futons having to be transported over large distances. During transportation, foldable futons are often stacked several high. A foldable futon located near the bottom of a stack maybe subjected to several thousands pounds of weight. This produces cracks and splits in the wooden planks that form the wooden futon frames, as well as bending in the rails that form the metal futon frames. The failure rate, however, of wooden futon frames, due to transportation, often far exceeds the failure rate of metal futon frames. This often reduces the savings benefit derived from transporting the lighter wooden futon frames.\nWhat is needed, therefore, is a foldable futon that is formed from wood and provides strength and durability comparable to a metal futon frame."}
{"text": "A typical bicycle wheel is formed completely of metal, and includes a cylindrical hub connected by numerous metal spokes to a rim, and an axle assembly that may include a brake lying within the hub. Improved performance of bicycle wheels has been obtained by using molded plastic wheels formed of very high strength plastic material. Design of the wheel for injection molding results in a different type of wheel, wherein a limited number of heavy duty spokes are preferably utilized. U.S. Pat. No. Des. 244,272 by Raudman and Hoffman show a wheel of this type. While high performance plastics have many advantages, one important disadvantage is that they are easily damaged by heat, such as that which can be generated by a bicycle brake when repeatedly applied. Also, the limited number of heavy duty spokes in a molded plastic wheel, can result in poor absorption of shocks when applied along the axis of one of the heavy duty spokes. The limited number of spokes also can have disadvantages in providing a large open space into which a child's foot may be inserted. A molded plastic vehicle wheel, especially one designed for use on bicycles, which helps minimize the transmission of shock between the rim of the wheel and the bicycle frame, and which minimizes the possibility of heat damage to the wheel by a brake assembly mounted in the hub of the wheel, would be of considerable value in the production of bicycles as well as other vehicles."}
{"text": "This invention relates, in general, to scoring tools, and, in particular, to scoring tools for wallpaper removal."}
{"text": "Recently, as an era of a high oil price advents, efforts to save fuel costs by improving energy efficiency of marine vessels and to secure an environmentally friendly feature in operation of a marine vessel are being made.\nIn general, in an operation of a marine vessel, most energy is consumed by a main engine for propelling the marine vessel, and about 25% of fuel consumed for the operation of the main engine is practically disposed of in air as an exhaust fume. Accordingly, a variety of devices for recovering part of waste heat by using the exhaust fume are actively being introduced.\nFIG. 1 is a schematic view illustrating a waste heat recovery device for a marine vessel of a related art. Referring to FIG. 1, in the waste heat recovery device for a marine vessel of a related art, a heat recovery device (boiler) 121 is installed at an exhaust pipe through which an exhaust fume is exhausted from an engine 110 of a marine vessel and then heat of the exhaust fume is recovered (heated under uniform pressure) so that high temperature steam is generated to be used as a variety of energy sources.\nHowever, in the waste heat recovery device for a marine vessel of a related art, since waste heat of an exhaust fume is recovered through only the heat recovery device 121 having a single structure, waste heat of an exhaust fume that is still in a high temperature state even after having passed through the heat recover device 121 is not recovered and discharged into the air so that energy is wasted.\nAlso, a heat recovery device for recovering heat from coolers to cool the heat generated from the engine itself is not provided.\nOn the other hand, an operation rate of the engine 110 varies during the operation of a marine vessel. For example, the engine 110 of a marine vessel may be operated with a full load for about ⅔ of the entire operation days and with a load lower than the full load for the rest of the operation days. As such, when the engine 110 is operated with a load lower than the full load, a relatively low amount of exhaust fume is generated compared to the exhaust fume generated when the engine 110 is operated with a full load.\nA change in the amount of an exhaust fume causes a change in heat capacity flowing in with the Rankine cycle, which has a severe influence on efficiency of the Rankine cycle. Accordingly, electricity may be stably generated when the above issue is considered in the development of a waste heat recovery device for a marine vessel which may drive a turbine by recovering a large amount of waste heat of an exhaust fume of an engine compared to the related art."}
{"text": "The present disclosure relates to mechanical machine parts and more particularly, but not exclusively, to a type of machine part commonly known as a “spring plunger” or a “ball and spring plunger.”\nMechanical devices known variously as “plungers,” “spring plungers,” “ball and spring plungers,” and “press-fit plungers” are used throughout industry in various applications. Such devices may be used as components of a larger apparatus such as a rotatable work surface to provide a locating feature which allows a user to reliably return the apparatus to a home position, or to provide a positioning feature which allows a user to reliably rotate the apparatus to one of a plurality of predetermined orientations, e.g., a 15 degree angle, a 30 degree angle, a 45 degree angle, etc. Related uses may include indexing, locating, and positioning applications in dies, special tooling, and automatic feed devices. Alternately, such devices may be used as components of an electrical circuit extending between the separable parts of an apparatus. In such uses, the device may be inserted into a conductive socket disposed in the transverse face of a first part, and oppose a conductive contact plate affixed to the transverse face of a second part. In such an arrangement, the device biases a conductive ball or nose against the contact plate, providing a durable electrical contact that is tolerant of reasonable variations in the distance separating the opposing faces of the respective parts. Related uses may include electrical supply or contact sensor applications in dies, automatic feed devices, and other machinery.\nTypical constructions of such devices include non-adjustable collared plungers and adjustable threaded plungers. Collared plungers may have a body that consists of a main barrel having an essentially constant outer diameter and a proximal collar surrounding a plunger opening. The collar typically extends laterally outward from the main barrel of the device, and serves to position the device over a blind aperture drilled or otherwise formed in the front or face of a receiving part. Collared plungers may be friction fit into such an aperture, but can loosen or separate from the receiving part in environments that experience substantial vibration and/or thermal cycling because the “tightness” of the friction fit is limited by the force that can be applied to the collar of the device during installation. Such devices are typically inserted into an aperture by lightly hammering or pressing the proximal end of the device, particularly the collar, to drive the device into the aperture, however, when excessive force is applied to the collar it may become damaged. Such damage may include distortion of the outer rim of the collar, which may prevent the collar from sitting flush against the surface of the receiving part, and distortion of the inner rim of the collar, which may prevent the plunger from returning to a fully extended position. Either type of damage may cause the assembly to fail to meet specifications. Typical uses of collared plungers are also restricted by the nature of their construction in that (1) the collar of the device typically rests above the front or face of the receiving part, which exposes the device to potentially damaging strikes from opposing cycling parts or other loose matter and (2) the device must be installed through the front or face of the receiving part, which restricts the range of applications or design configurations that may make use of such a device without substantial design modifications.\nThreaded plungers overcome many of the drawbacks associated with collared plungers, but require significant machining and installation preparation. Threaded plungers may have a body consisting of a threaded shaft and a proximal slotted rim surrounding a plunger opening. Such plungers may also include a plastic locking element and/or second slot disposed across a distal wall of the device. Threaded plungers are typically screwed into an aperture that has been tapped to provide complementary threading, and the plastic locking element, if present, serves to frictionally engage the aperture threads to increase the torque that must be applied to install, adjust, or remove the device. Such devices are typically installed by using a plunger wrench that engages the slotted rim of the device while minimizing contact with the plunger, permitting one to screw the device into a threaded aperture in the front or face of a receiving part. Such devices may also be installed by using a narrow-blade screwdriver that engages the second distal slot, if present, permitting one to screw the device into a threaded aperture through the rear of the receiving part. If the device is installed through the front or face of a receiving part with a screwdriver rather than a plunger wrench, the blade of the screwdriver will depress the plunger of the device past the outer surface of the slotted rim during installation. Such excessive travel may overcompress an internal spring, altering the plunger force profile and damaging the device.\nThus there is also a need for a plunger device that may be installed into a drilled or formed aperture through either the front or the rear of a receiving part without additional preparation of the aperture, such as tapping. There is also a need for a non-threaded plunger device that is better able to resist vibration and thermal cycling."}
{"text": "Capillary electrophoresis (CE) is an electrophoresis technique utilizing small bore capillaries. CE provides methods for the separation of ionic species including macromolecules. The efficiency of CE can be relatively high, i.e. in excess of 400,000 theoretical plates, and thus is being explored for a number of different applications.\nA typical CE system includes a 50-100 micrometer internal diameter silica capillary tube filled with a suitable electrically conducting buffer. The outlet end of the capillary is immersed in a reservoir containing the buffer and an electrode. A sample containing ions of interest is introduced into the inlet end of the capillary and then the inlet is placed in another reservoir containing the buffer and another electrode. Since a small diameter capillary is used in CE, a relatively high applied voltage can be used without the generation of thermal gradients in the capillary; thus the electrodes are connected to a power supply capable of delivering .sup..about. 30 kV per 100 cm of capillary. A detector is placed between the two electrodes to permit detection of various ionic species migrating in the capillary. A detector so positioned is often referred to as an on-column detector. Typically, an integrator is attached to the detector, such that the peak areas may be measured.\nThe movement of the sample ions of interest is controlled by two factors: the electrophoretic velocity and the electroosmotic flow velocity. The total migration velocity is the vector sum of these two terms.\nElectrophoretic migration is the migration of the sample ion towards the oppositely charged electrode under the influence of the electric field. The electrophoretic mobility of any particular ion is the electrophoretic velocity per unit field strength.\nElectroosmotic flow (EOF) is the bulk flow of the buffer in the capillary. EOF is due to the charge of the inside surface of the capillary which is in contact with buffer containing mobile counterions. For example, an unmodified silica capillary surface comprises silanol (Si--OH) groups that are negatively charged (Si--O.sup.-) when the pH of the buffer is greater than about 2, and positively charged (Si--OH.sub.2 +) when the pH is less than about 2. Alternatively, hydrophobic cations may be adsorbed onto the inside surface of the capillary to obtain a positively charged surface at higher pHs.\nWhen the surface is negatively charged, then the mobile counterions, for example, sodium ions (Na.sup.+), migrate under the influence of the electric field and in the process drag the bulk solvent with them. Thus the direction of the electroosmotic flow is from the positive to the negative electrode when the surface is negatively charged.\nWhen the surface is positively charged, then the mobile counterions of the positively charged surface, e.g. biphosphate ions (HPO.sub.4.sup.2-), migrate under the influence of the electric field and in the process drag the bulk solvent with them. Thus the direction of the electroosmotic flow is from the negative to the positive electrode when the surface is positively charged.\nWhen the surface is not charged, then there is no electroosmotic flow, and any movement of analyte ions is due solely to electrophoretic mobility.\nThus, depending on the charge of the ions of interest, the nature and the extent of capillary surface charging and the polarity of the applied voltage, electroosmosis can augment, counteract or even override the electrophoretic movement. Since sample components to be determined must travel from the inlet end of the capillary to the detector which is located near the outlet end of the capillary, it is essential that they move in the desired direction. However, since the total migration velocity of the sample is the vector sum of the electrophoretic velocity and the electroosmotic flow velocity, it is possible that the charge of the sample is such that it would move away from the outlet electrode in the absence of electroosmotic flow; under these conditions the electroosmotic flow velocity of the bulk solution must be greater than the electrophoretic mobility of the analyte.\nThe detector used in the CE system is very important, and the type of detector used will usually depend on the properties of the compounds under analysis. Currently, there are a number of different detector schemes utilized in CE. These include direct and indirect photometric detection, direct and indirect fluorescence detection, as well as suppressed and non-suppressed conductometric detection. Other types of detection which may be utilized are mass spectrometry, radiometry and other electrochemical methods such as amperometry. These methods may be used either on-column or end-column, i.e. at some point after the outlet electrode.\nWhile most organic molecules of interest display significant ultraviolet absorption such that direct photometric detection is practical, this is not the case for many inorganic ions or aliphatic carboxylic acids that display very low optical absorption. Thus indirect photometric or fluorescence detection may be utilized in these cases. For example, indirect fluorometric detection is described in Gross et al., Anal. Chem. 62:427-431 (1990); Bachmann et al., Journal of Chrom. 626:259-265 (1992); and Gross et al., Journal of Chrom. 480:169-178 (1989). Indirect photometric detection is described in Foret et al., Journal of Chrom. 470:299-308 (1989); Foret et al., Electrophoresis 11:780-783 (1990); and Henshall et al., Journal of Chrom. 608:413-419 (1992). Additionally, since electrical mobility is an intrinsic property of all ions, detection based on conductivity can be a desirable method for many uses of CE.\nHowever, conductivity detectors are nonselective bulk property detectors. The signal arises from the difference in equivalent conductance or mobility of the charge carrier electrolyte ion and the analyte ion. In CE, a large difference in mobility of the carrier electrolyte ion and the analyte ion leads to excessive peak tailing/fronting, which means that there are practical restrictions on the choice of the eluent ion. This conflict between optimum sensitivity and separation efficiency represents the ultimate limitation of nonsuppressed conductivity detection in CE.\nElectrolyte suppression, or the post run alteration of the electrolyte buffer such that the background \"noise\" of the buffer is decreased, has been explored for ion chromatography (see for example U.S. Pat. Nos. 3,897,213; 3,920,397; 3,925,019; 3,956,559; 4,474,664; 4,751,004; 4,459,357 and 4,999,098), and recently for CE (U.S. Ser. Nos. 07/771,336 and 07/771,597, filed Oct. 4, 1991, herein incorporated by reference). These systems are referred to as \"suppressed\" systems.\nWhen conductometric detection is utilized, aspects of the electropherogram from a CE run are different from an optically detected CE run. These phenomena are the subject of the present invention.\nOne problem with CE is that the sample volume introduced into the capillary may vary; samples are not easily injected onto the column by a fixed volume valve. This can be a problem in other chromatography systems as well, although many systems do employ a fixed volume valve. Instead, in CE, the samples are introduced into the capillary in several ways. Typical injection modes for CE are pressurization of the sample- or standard-containing vial for a fixed length of time (pressure injection), or the application of an electrical field for a fixed length of time (electrostatic injection). Thus CE sample injection methods are based on time of injection. This means that small variabilities in time, sample viscosity, pressure or hydrostatic height may result in variabilities of the sample volume injected. This may have a profound impact on the migration time and the size of the peaks.\nFurthermore, these potential variabilities may make quantification of sample peaks difficult. Quantification requires that the injected volume of the sample and the injected volume of the standard be essentially exactly the same, or that the volume of each be precisely known to compensate for injected volume differences.\nAnother problem with CE and other systems relying on electroosmotic flow is that it is difficult to determine the electroosmotic flow velocity, and to detect any variations in this rate. For example, in CE, the electroosmotic flow velocity in a bare silica capillary can be affected by material adsorbed on the wall from the previous injection, which may alter the flow rate. Similarly, the flow rate of different capillaries may be different due to a variety of factors. The electroosmotic flow velocity will also change as a result of a change in applied voltage. Since the flow velocity will affect the migration times of the analyte peaks, i.e. total migration velocity of the analyte, this may be a serious problem. This takes on an increased significance if an integrator is used, since the integration is dependent on time; thus a decrease in the flow rate for a sample run may result in an increase in the peak area, and vice versa.\nThese two limitations are addressed by the present invention."}
{"text": "1. Field of the Invention\nThis invention relates in general to enterprise resource planning systems performed by computers, and in particular, to a method and apparatus for accessing an enterprise resource planning application over the Internet using Java.\n2. Description of Related Art\nWith the fast growing popularity of the Internet and the World Wide Web (also known as “WWW” or the “Web”), there is also a fast growing demand for Web access to Enterprise Resource Planning (ERP). However, it is especially difficult to use ERP software with the Web. One of the problems with using ERP software on the Web is the lack of correspondence between the protocols used to communicate in the Web with the protocols used to communicate with ERP software.\nFor example, the Web operates using the HyperText Transfer Protocol (HTTP), the eXtensible Markup Language (XML), the eXtensible Style Language (XSL) and the HyperText Markup Language (HTML). The use of this protocol and languages result in the communication and display of graphical information that incorporates hyperlinks. Hyperlinks are network addresses that are embedded in a word, phrase, icon or picture that are activated when the user selects a highlighted item displayed in the graphical information.\nHTTP is the protocol used by Web clients and Web servers to communicate between themselves using these hyperlinks. HTML is the language used by Web servers to create and connect together documents that contain these hyperlinks. The HTML syntax and commands are specified by the web browsers, and cannot be extended by users.\nXML and XSL are anticipated to be the next generation of web languages. XML is the language used by Web servers to create and connect together documents that contain user defined structures. XSL is the language used by web browsers to convert XML documents into HTML for the purposes of display. The validity of an XML document is defined by a Document Type Definition (DTD), which an XML parser uses to ensure that an XML document is valid.\nWeb servers are extensible via a number of APIs. The Common Gateway Interface (CGI) is a standard interface for executing programs external to the web server. Java is a programming language and environment defined by Sun Microsystems. Java Servlets are a widely available interface for executing Java programs within the web server. Java Server Pages (JSP) is a web page interface for specifying Java commands that are to be executed in the web server. Visual Basic Script (VBScript) is a programming language defined by Microsoft. Active Server Pages (ASP) is a web page interface specified by Microsoft for specifying VBScript commands that are to be executed in the web server.\nIn contrast, most ERP software provides an application programming interface (API) for accessing ERP functionality from external programs. Each vendor of an external program may have its own implementation of the API, for example, Peoplesoft has an API, called the Message Agent API.\nThus, there is a need in the art for methods of accessing Enterprise Resource Planning software across the Internet network, and especially via the World Wide Web. Further, there is a need for simplified development environments for such systems."}
{"text": "1. Field\nEmbodiments of the present inventive concept relate to a display device capable of facilitating wiring processes between a light source module and a power supply configured to apply driving signals to the light source module.\n2. Description of the Related Art\nA liquid crystal display (LCD) device is a type of flat panel displays (FPDs), which is most widely used these days. An LCD device includes two substrates including electrodes formed thereon and a liquid crystal layer interposed therebetween. Upon applying voltage to two electrodes, liquid crystal molecules of the liquid crystal layer are rearranged, thereby adjusting an amount of transmitted light.\nSuch an LCD device may include a backlight unit. The backlight units may be categorized into a direct type, an edge type, a corner type, and the like according to a position of the light source.\nIn recent years, a light emitting diode (LED), which is characterized by low power consumption and high efficiency, is usually used as a light source of a backlight unit.\nFurther, a light source module including a plurality of LEDs mounted on a printed circuit board (PCB) may be used for the backlight unit.\nThe light source module may receive driving signals from an external power supply and may emit light. For this purpose, wires may be disposed between the light source module and the power supply to connect therebetween.\nThe direct-type backlight unit may include a large number of light source modules, and thus more wires may be disposed between the light source modules and the power supply.\nIn particular, a curved-type display device having such a direct-type backlight unit may include, compared to a general display device, more light source modules. This is mainly because the curved-type display devices are mostly provided in a quadrangular shape and curved along the longitudinal direction of a long side of the quadrangular shape, which is described below in detail.\nThat is, when the light source modules are arranged along the longitudinal direction of a short side of the curved-type display device, it is advantageous in that less light source modules are required. In this case, however, the light source modules may be curved, which incurs unstable operation of the light source modules. Accordingly, light source modules used for conventional curved-type display devices are arranged along the longitudinal direction of the long side of the display device. In such cases, the longitudinal direction of each light source module may correspond to the longitudinal direction of the short side which is not curved, such that the light source modules may not be curved. However, due to such structure, the conventional curved-type display device may require more light source modules. Accordingly, in the conventional curved-type display device, wirings between the light source modules and the power supply may become disadvantageously complicated.\nIt is to be understood that this background of the technology section is intended to provide useful background for understanding the technology and as such disclosed herein, the technology background section may include ideas, concepts or recognitions that were not part of what was known or appreciated by those skilled in the pertinent art prior to a corresponding effective filing date of subject matter disclosed herein."}
{"text": "The present invention is directed to a phase-control power controller that supplies a specified power to a load, and more particularly to a voltage converter for a lamp that converts line voltage to a voltage suitable for lamp operation.\nSome loads, such as lamps, operate at a voltage lower than a line (or mains) voltage of, for example, 120V or 220V, and for such loads a voltage converter that converts line voltage to a lower operating voltage must be provided. The power supplied to the load may be controlled with a phase-control power circuit that typically includes an RC circuit. Moreover, some loads operate most efficiently when the power is constant (or substantially so). However, line voltage variations are magnified by these phase-control circuits due to their inherent properties (as will be explained below) and the phase-control circuit is desirably modified to provide a (nearly) constant RMS load voltage.\nWhen the phase-control power controller is used in a voltage converter of a lamp, the voltage converter may be provided in a fixture to which the lamp is connected or within the lamp itself. U.S. Pat. No. 3,869,631 is an example of the latter, in which a diode is provided in the lamp base for clipping the line voltage to reduce RMS load voltage at the light emitting element. U.S. Pat. No. 6,445,133 is another example of the latter, in which transformer circuits are provided in the lamp base for reducing the load voltage at the light emitting element.\nFactors to be considered when designing a voltage converter that is to be located within a lamp include the sizes of the lamp and voltage converter, costs of materials and production, production of a potentially harmful DC load on a source of power for installations of multiple lamps, and the operating temperature of the lamp and an effect of the operating temperature on a structure and operation of the voltage converter."}
{"text": "1. Field\nExemplary embodiments of the present invention relate to a backlight unit and a display apparatus having the backlight assembly and, more particularly, exemplary embodiments of the present invention relate to a backlight unit providing a liquid crystal display (LCD) apparatus with light and an LCD apparatus having the backlight assembly.\n2. Discussion of the Background\nIn general, an LCD apparatus includes an LCD panel to display an image and a backlight unit disposed under the LCD panel to provide the LCD panel with light.\nThe LCD panel includes a first substrate having a thin film transistor (TFT) and a pixel electrode electrically connected to the TFT, a second substrate having a color filter corresponding to the pixel electrode and a common electrode, and a liquid crystal layer disposed between the first substrate and the second substrate.\nThe backlight unit includes a receiving container, a plurality of lamps disposed in parallel in the receiving container, and an optical member disposed over the lamps to improve optical characteristics of the light generated by the lamps. The optical member includes a light-diffusing plate diffusing the light to improve luminance uniformity.\nThe optical member is in general spaced apart from the lamps to improve diffusion of light. As a result, the thickness of the backlight unit increases."}
{"text": "Nonwoven fabrics have been used for a wide variety of applications for at least fifty years. Nonwoven fabrics are textile-like materials produced directly from a web of fibers so as to eliminate the many time consuming steps involved in converting staple length fibers into woven or knitted goods. In one method of making a nonwoven fabric, a web of fibers is produced, e.g. by carding or air laying techniques, and the fibrous web is then strengthened by the application thereto of a polymeric binding agent. In another method of making a nonwoven, the fibrous web is subjected to fluid forces which serve to entangle the fibers, thus providing strength to the finished material. Nonwoven fabrics are inherently porous structures, i.e. they comprise openings or pores allowing for the passage of fluids such as air and water or aqueous solutions. In addition, nonwoven fabrics may be tailored so as to have good softness, drapeability, and tactile impression. Owing to their desirable characteristics, nonwovens have been employed as facing materials for absorbent products such as disposable diapers, sanitary napkins, incontinent devices, wound dressings and the like.\nMore recently, efforts have been made to produce porous or liquid-permeable facing materials for absorbent products by using plastic films as the starting materials. For example, it is known to produce an apertured plastic film by placing a heated thermoplastic sheet material on a patterned perforated surface and applying a vacuum thereto. The vacuum pulls the softened sheet material through the perforations, thereby causing the film to rupture and form apertures.\nU.S. Pat. No. 3,929,135 to Thompson et al discloses perforated topsheet materials for absorptive devices such as sanitary napkins, incontinent pads, bandages and the like. These topsheet materials are constructed from liquid impervious materials such as low density polyethylene and comprise a plurality of tapered capillaries each of which has a base opening in the plane of the top sheet and an apex opening which is remote from the plane of the top sheet. The tapered capillaries disclosed by Thompson et al. are preferably provided in the form of a frustum of a conical surface and have an angle of taper of from about 10.degree. to 60.degree..\nU.S. Pat. No. 4,324,276 to Mullane discloses an apertured formed film having a caliper of less than about 0.030 inch (0.075 cm), an open area of at least 35% and a plurality of apertures at least 75% of which have an equivalent hydraulic diameter (EHD) of at least 0.025 inch (0.064 cm). The apertured formed film disclosed by Mullane et al. is useful as a topsheet for disposable absorbent products of the type mentioned above.\nU.S. Pat. No. 4,839,216 to Curro et al. discloses a debossed and perforated plastic material produced by providing a starting film on a perforated forming surface and applying an unconstrained liquid stream to the upper surface of the starting film. The liquid stream has sufficient force and mass flux to cause the film to be deformed toward the forming surface, such that the material acquires a substantial three-dimensional conformation, and to cause perforations to be created therein.\nEuropean Patent Application 0 304 617 in the name of Kao Corporation discloses a covering sheet for a sanitary article. The covering sheet comprises an opaque, hydrophobic film having land portions and recesses, said recesses being formed to have a bottom portion and side walls. The side walls have a slanting part which is provided with an opening such that the slanting part is not covered by the land portion. This patent states that the opening is always exposed to sight when it is looked down at.\nU.S. Pat. No. 4,690,679 discloses an apertured film comprising a first layer of a first polymeric film and a second layer of a second polymeric film. Apertured films wherein the apertures have average equivalent circular diameters ranging from about 0.010 inches (0.0254 cm) to about 0.030 inches (0.0762 cm) are disclosed as being useful as covering materials for absorbent products.\nOther patents relating to apertured films and methods and apparatus for making the same include U.S. Pat. No. 3,632,269 to Doviak et al. and U.S. Pat. No. 4,381,326 to Kelley."}
{"text": "1. Field of the Invention\nThe present invention relates to a monoazo compound used to color petroleum products and hydrocarbon solvents and for other purposes, and a dense solution of said monoazo compound in a hydrocarbon solvent.\n2. Description of the Prior Art\nAs coloring agents for synthetic resins, lacquers, natural and synthetic waxes, solid paraffins, etc., various monoazo oil-soluble dyes listed in the Color Index (e.g., C.I. Solvent Colors: YELLOW 2, 7, 9, 10, 11, 12, 14, 16, 18, 56, 58; ORANGE 2, 7; RED 2, 24) are known.\nIndustrial organic solvents and petroleum products such as gasoline, lamp oil and light oil are colored for the purpose of their identification. For example, fuel oil is colored with an oil-soluble dye with improved solubility in hydrocarbon organic solvents. In this case, the oil-soluble dye is usually used as a solution of the maximum possible concentration in the hydrocarbon organic solvent.\nProposals concerning oil-soluble azo dyes for these uses are disclosed in the specification for British Patent No. 1142239, Japanese Patent Examined Publication Nos. 17390/1981 and 36940/1982 and Japanese Patent Publication Open to Public Inspection (hereinafter referred to as Japanese Patent O.P.I. Publication) Nos. 2525/1971, 94029/1975, 99959/1980 and 187456/1983.\nFor example, the specification for British Patent No. 1142239 describes a monoazo dye synthesized from m-toluidine and N,N-diethylaniline or N,N-dipropylaniline, and a dense solution thereof used to color petroleum and for other purposes.\nAlso, Japanese Patent O.P.I. Publication No. 2525/1971 describes that the dense solution prepared by dissolving a monoazo dye represented by the following formula: ##STR2## in an aromatic hydrocarbon solvent such as xylene is capable of imparting a bluish red color on petroleum products.\nJapanese Patent O.P.I. Publication No. 99959/1980 discloses a monoazo dye capable of forming a salt with aliphatic amine obtained by reaction of an aniline substituted by a long-chain alkyl having 8 to 18 carbon atoms and a pyrazolone derivative having a carboxyl group or .beta.-naphthol.\nJapanese Patent O.P.I. Publication No. 94029/1975 discloses a monoazo dye obtained by reaction with, for example, a alkylamino-substituted .beta.-naphthol derivative having 4 to 21 carbon atoms as a coupling component, stating that this red dye dissolves well or limitlessly in aromatic and aliphatic substances, and that it can be used to color automobile fuels, varnish, wax and fat.\nThe above-mentioned oil-soluble dyes listed in the Color Index have any one or more of the following drawbacks:\n1) Solubility in hydrocarbon solvents, particularly paraffin solvents, is low. The saturation concentration of this kind of conventional oil-soluble dyes in paraffin solvents at normal temperature does not exceed 15% by weight. For example, the saturation concentration in hexane at normal temperature is not higher than 5% by weight.\n2) Coloring power for petroleum products such as fuel oil is weak. Colored fuel oil should make it possible to distinguish between petroleum products to be taxed by coloring with coloring agents of concentrations of the order of some dozens of ppm.\n3) Because of high crystallinity and high intermolecular aggregating force, the oil-soluble dye is easy to extract from the colored system. A countermeasure is therefore necessary against tax evasion by discoloring the petroleum products to be taxed.\n4) In the case of dyes of long-chain alkyl-substituted structure, the high molecular weight lowers the coloring density per unit part by weight.\nAlso, the dye described in Japanese Patent O.P.I. Publication No. 99959/1980 poses a problem in coloring power, whose solubility in ligroine or kerosine is 0.5 to 1.5 g/10 cc.\nThe object of the present invention is to provide a monoazo compound which is excellent in compatibility with, solubility in and coloring power for almost all water-immiscible organic liquids, such as various hydrocarbon solvents, petroleum solvents, fuel oils, lubricating oils, plastics and hydrophobic synthetic fibers, and which is suitable for use to color molten plastics and dope dyeings and to color thermal transfer ink because of the low melting point, and a dense solution thereof which is free of the problems of powder formation and dust scattering because of the liquid form, which can easily be prepared by simple dissolution operation, which quickly and uniformly diffuses and dissolves in coloring subjects such as petroleum products, and which is capable of coloring fuel oil etc. to a desired distinguishable color density in small amounts."}
{"text": "1. Field of the Invention\nThe present invention relates in general to an automotive selected gear sliding type manual transmission, and more particularly to a gear operating mechanism of the transmission, which has means for positively preventing the undesirable gear disengagement phenomenon.\n2. Description of the Prior Art\nIn a manual transmission of a selected gear sliding type, a so-called \"gear disengagement\" sometimes happens particularly in a long-used transmission. The gear disengagement is an undesirable phenomenon in which the gear assuming an operating or engaged position suddenly and unexpectedly disengages from the associated gears thereby suddenly breaking the power transmitting connection between the engine and the driven road wheels. In order to solve such undesirable and dangerous phenomenon, various measures have been hitherto proposed. However, some of them have been made at the sacrifice of handling comfortableness and production cost of the transmission."}
{"text": "Cardiovascular disorders come first in the world as far as lethality and invalidization of people are concerned.\nIn the recent years there has been an awakening to the need for a more systematized approach to the problem of physical exercise. Particularly, individuals in the over-thirty years age bracket are having routine exercise prescribed by their physicians as a preventive measure against heart desease and other associated cardiovascular disorders. Numerous recent publications have recognized a correlation between physical fitness and exercise and have delineated detailed regimens or programs of physical exercise intended to yield the desired results. This physical exercise must be performed only under medical supervision with constant checks on physical loads and the physical activity of the exerciser.\nBesides, in the course of medical examination it is often necessary to know about a person's heart activity when the person is subjected to various physical loads."}
{"text": "An arrangement is already known from DD 260 847 A3 designed to provide smooth, linear elastic support which offers lateral stability and is torsion-proof. The arrangement comprises pre-formed leaf springs which are slit and which are firmly fixed at the ends of the resulting spring tongues into the interlocking spring mechanism of the components. Here, the different pre-formed leaf springs lie parallel to each other, the arrangement having a symmetrical form. The fixing points of the pre-formed leaf springs concerned lie at a level which is parallel to the direction of spring movement, the total width of the spring tongues, which are each fixed within the movable components, being the same for a pre-formed leaf spring and at the same distance from the fixing point. Due to the angle position of the leaf springs concerned during bending, a traverse movement of the actuator is generated which runs at right-angles to the direction of movement.\nWO 97/13603 according to the FIG. 1, together with the related description in the penultimate paragraph on page 9, describes a processing machine for manufacturing optical surfaces and formed shells according to individual data specifications in a single processing stage. Here, a tool holder which oscillates in the Z direction is fitted, which is in turn assigned a workpiece holder with a spindle on a slide which can be moved in the Z axis. Here, the slide is brought into a position for preadjustment, and is then clamped firmly in place for the duration of the processing. A dynamic drive according to the invention, which is moved during the processing sequence, is not described in this publication.\nAn object of the present invention is to form and arrange a lens processing machine and the related linear drive so as to enable simple, low-cost processing of the workpiece.\nThe present invention provides a highly dynamic lens processing machine that includes a workpiece holder having at least one dynamic linear drive in the direction of a second oscillation axis O2 which is arranged in parallel to the first oscillation axis O1. Furthermore, a runner is guided over a leaf spring bearing arrangement formed from leaf springs in the housing in a coaxial direction to the oscillation axis O1, the leaf spring bearing arrangement being formed in two parts, and comprising a guidance arrangement and a coupling arrangement attached to said guidance arrangement. Here, the linear drive can take the form of a linear motor, or equally a slide-spindle pair. The drive can be regarded as being dynamic when the processing or milling sequence is actively supported or influenced. Here, an incremental, constant or step-by-step trigger is provided.\nThis enables the oscillation axis 01, 02 required for processing the workpiece or lens to be partially arranged in the tool holder, and in addition to this, partially in the workpiece holder. As a result, in particular the required strokes of the oscillation axis O1, O2 concerned can be used for processing on the one hand, as well as the required processing or oscillation frequencies, the oscillation of an axis or a combination of the oscillations of both axes. In this regard, it is also advantageous that the relationship of the oscillation stroke of the tool holder to the oscillation stroke of the workpiece holder lies between 0.08 and 1, in particular 0.13. When manufacturing a lens, the lens surface, which requires highly dynamic triggering for the cutting chisel, can be guaranteed due to the highly dynamic first oscillation axis O1 of the leaf spring bearing arrangement. The milling processing of the edge is conducted using the dynamic second oscillation axis O2, which comprises a significantly larger stroke than the first oscillation axis O1.\nAs a result, a simple and, in this respect, maintenance-free linear drive for the tool holder and/or the workpiece holder is guaranteed. Due to the use of a leaf spring bearing arrangement, the use of expensive slide and/or rolling bearings is not required. The stroke of the oscillation axis concerned, which is reduced in this respect due to the use of the leaf spring bearing arrangement can be compensated by overlaying or combining the aforementioned double oscillation arrangement. As well as the oscillation frequency of between 25 and 100 Hz, an oscillation path of between 1 and 20 mm is preferably provided in order to process the optical surface of a lens blank.\nFor this purpose, it is also advantageous that a second tool holder which is arranged in a stationary position is provided, towards which the workpiece can be moved via the second oscillation axis O2. As well as the overlay or combination of the two oscillation axes, the oscillation axis O2 which is assigned to the workpiece holder can also be used in order to process the workpiece. In this case, an oscillation of the tool holder is not necessary, so that said tool holder is arranged in a stationary position, for example in the form of a milling cutter, and the workpiece conducts the oscillation. When a rotating chisel is used for processing the plastic lens blanks, the first oscillation axis O1, i.e. the dynamic tool drive, is preferably used in order to process the optical surface, the second oscillation axis O2 of the workpiece holder being used to process the edge or to process the phases, and for the basic roughing-down work. Here, the use of a milling cutter is specified as an alternative to the rotating chisel.\nFurthermore, it is advantageous that the leaf spring bearing arrangement comprises a coupling arrangement and a guidance arrangement, the coupling arrangement and/or the guidance arrangement being indirectly or directly connected with the runner and/or with the housing. The use of a coupling and guidance arrangement in each case, which in this respect forms a structural unit, guarantees the necessary rigidity of the bearing arrangement on the one hand, and the optimisation of the vibration behaviour of the overall arrangement on the other, in particular taking into account the frequencies of between 25 and 100 Hz to be specified here.\nIt is also advantageous that the guidance arrangement is formed from a first guide segment and a second guide segment, which is arranged at a distance in the direction of the oscillation axis O1. This guarantees double, front bearing support for the runner.\nFor this purpose, it is advantageous that the guide segment concerned comprises at least one bar-formed or plate-formed reinforcement element, and that the reinforcement element concerned is preferably carbon fibre-strengthened, or is formed of carbon fibre, the reinforcement element concerned being adhered to an adapter element which is preferably welded to the guide segment. Due to the use of reinforcement elements, the elastic part of the guide element concerned is reduced on its end sections and in the central section, i.e. the section where the runner or housing attachment is located. Here, the elastic length is approximately reduced to between 5 mm and 15 mm. Its structure made of, or using, carbon fibre guarantees the best possible weight ratio. Using the adapter element, the adhesive attachment next to the support surface can also cover the edge areas or front face sides of the adapter element.\nHere, it is specified in an advantageous manner that the runner in the direction of the oscillation axis O1 comprises a front face side and a rear face side, the guide segment concerned being affixed to the runner in the area of the face side in order to form a first and a second runner bearing. The use of two guide segments results in the runner being retained on both sides on its two face sides. Due to the face side arrangement of the guide segments on the runner, the interim area which is thus created can be used to attach further runner bearing parts.\nOf particular significance to the present invention is the fact that the external end sections of the guide segments in the radial direction to the oscillation axis O2 are connected to the coupling arrangement. The connection between the guide segments to the coupling arrangement guarantees the necessary coupling or restraint, and thus the required discontinuation of the vibration behaviour of the two guide segments. The external ends of each of the guide segments concerned act here as an optimum attachment point in order to attach the coupling arrangement.\nIn relation to the structure and arrangement according to the invention, it is advantageous that the coupling arrangement is formed from a first coupling segment and a second coupling segment, which are arranged opposite to each other in relation to the oscillation axis O1. This enables the external ends of the guide segments to be connected, which are also aligned in a radial direction to the oscillation axis, and which are also arranged opposite in relation to the oscillation axis.\nFurthermore, it is advantageous that at least one end of the first guide segment and at least one end of the second guide segment are connected via at least one coupling segment. This mechanical coupling of the guide segments via the coupling segment concerned guarantees an optimum vibration behaviour on the guide segments, in particular taking into account the different natural frequencies of said guide segments. The natural frequencies are influenced by the coupling and guidance arrangement, and are thus moved into a frequency range outside that of the processing frequency.\nHere it is advantageous that the connection between the guide segment and the coupling segment takes the form of a plug connection. The guide segment concerned comprises several recesses at each end, into which the corresponding pins of the coupling segment can be inserted. A weld connection can also be provided in addition to this plug connection.\nIt is also advantageous that the guide segment has a rhomboidal form or an oval double form, and that it is affixed to the housing via a housing bearing. Due to this basic form, the attachment to the runner or runner bearing is also guaranteed at the height of the runner, as well as the attachment point on the housing, i.e. the housing bearing, whereby starting from this central connection located in the area of the runner, the ends which extend outwards in a radial direction allow the continued attachment of the guide segment to the coupling segment concerned.\nFurthermore, it is advantageous that a carrier element, which takes the form of a steering rod, is provided between the housing and the runner in a radial direction to the oscillation axis O1, whereby the steering rod is connected with the housing via a first joint and with the runner via a second joint. The steering rod therefore represents the core part of the leaf spring bearing arrangement, in particular taking into account the vibration behaviour of said bearing arrangement which occurs, and the necessary guidance of the movements which are created as a result. Here, the connection rod is triggered by oscillating movements of the runner, and thus completes a horizontal swing around the first joint on the side of the housing. Here, the joints are preferably formed using elastic spring joints, which do not comprise any slackness, and which have an adequate degree of rigidity alongside their degree of freedom.\nHere, it is advantageous that at least one coupling segment is connected with the steering rod via a third joint. The attachment of the coupling segment to the steering rod guarantees an oscillation movement of the coupling segment which depends on the horizontal swing movement of the steering rod, which in turn is transferred to the relevant ends of the guide segments to which it is connected. Here, the third joint is also formed as an elastic flexible joint, so that the coupling segment is subjected to the minimum bending stress when the steering rod is swinging.\nFinally, it is advantageous that the joint is formed as a rolling bearing or as an elastic flector. The formation of the joint as an elastic flector or spring joint guarantees the absolute zero flexibility of the attachment or connection of the individual segments, which, taking into account the desired processing frequency of between 25 and 100 Hz, and the masses to be moved, guarantees adequate rigidity on the one hand, and sufficiently low maintenance on the other, together with the ensuing necessary load level.\nFor this purpose, it is also advantageous that a distance A is provided in a radial direction to the oscillation axis O1 between the first joint and the second joint, and a distance B is provided between the first joint and the third joint, the ratio of A to B of the two distances being 2. An aspect ratio of 2 guarantees that the excursion of the coupling segment in the direction of the oscillation axis is exactly half the amount as the excursion of the runner itself. Any movement of the runner towards the oscillation axis, which brings with it a contradictory alteration in the opening angle concerned, α1, α2, of the rhomboid-formed guide segments, causes the ends of the guide segments to be moved towards the oscillation axis. This movement of the ends corresponds to the movement of the coupling segment which is arranged in the centre of the steering rod. The variation of this distance ratio should be specified here according to the structure of the guide segments.\nFinally, it is advantageous that the second joint also guarantees a relative movement between the steering rod and the runner in a right-angled direction to the oscillation axis O1. Due to the tipping motion of the steering rod, the distance between the first and second bearing or bending point is altered, so that the necessary longitudinal compensation at right angles to the oscillation axis, or in the direction of the steering rod, is guaranteed by the second joint.\nFor this purpose, it is also advantageous that at least one coupling segment is connected with the third joint via a coupling member. The use of the coupling member makes it possible to vary the position of the coupling segment relative to the steering rod independently of the aforementioned attachment, while at the same time, taking into account said attachment.\nThe coupling element is thus arranged at a distance from the bearing point or swivelling axis of the third joint, whereby using the coupling member, the displacement path which is formed by the steering rod is guaranteed in accordance with the position of the third joint on the steering rod. Here, several coupling segments, in particular those which are arranged adjacent to a steering rod, can also be connected to the steering rod via a coupling member.\nFurthermore, it is advantageous that the coupling segment and/or the coupling member has a U-form, L-form or triangular transverse profile. This structure ensures the necessary rigidity, taking into account bending and torsion stresses during operation.\nHere, it is advantageous that the coupling segment comprises a connection point with the ends of the guide segments which takes the form of an elastic joint. The elastic joint has a flat structure, so that the relative excursion movement between the coupling segment and the guide segment concerned can be compensated. The aforementioned rigidity is therefore cancelled in the end sections.\nIt is also advantageous that the distance between the housing and the steering rod and/or the first joint can be adjusted using an adjusting device. Taking into account the required distance ratios between the firs and second bearing on the one hand, and the second and third bearing on the other, the adjustment device guarantees that the required distance ratio will be set, regardless of any given processing tolerances. The distance between the first bearing and the housing, or between the steering rod and the housing, can be created variably, so that the distance between the first and second bearing can be varied, depending on the distance produced between the second and third bearing.\nIt is also advantageous that the runner runs on bearings on the housing via the steering rod and the guide segment, the guide segment runs on bearings on the coupling segment, and the coupling segment runs on bearings on the housing via the steering rod. The frictional connection circuit thus created between the coupling and guidance arrangement on the one hand, and the housing and runner on the other, guarantees an optimum natural vibration frequency of the runner and an adequate displacement of the natural vibration frequency of the two arrangements in an area outside that of the processing frequency. When a cutting chisel is used as a tool and when plastic lenses are being processed, the cutting force to be applied to the chisel approximately totals between 50 and 200 Newton. Taking into account the processing frequency of between 25 and 100 Hz in relation to the oscillation frequency, and taking into account the masses of this linear drive required here, i.e. the runner, the tool holder and the motor actuator, it becomes clear, that the resulting vibration behaviour of the drive due to the required processing precision must be at an optimum level. Here, it is advantageous that the tool holder comprises a height adjustment feature, so that taking into account the alignment of the tool and the workpiece next to the processing axes, i.e. of the oscillation axis on the one hand, and a processing axis arranged radially to the workpiece on the other, the necessary adjustment to a spatial axis which is arranged at right angles to these two axes is guaranteed.\nIt is furthermore advantageous that at least one U-formed linear motor is assigned to the runner, which motor consists of the stator, which is formed from two magnetic tracks arranged in parallel, the magnetic tracks retaining the actuator between them. The motor coil of the actuator or the stator is air-core, poured in plastic and therefore very light. The forces of attraction of the opposite magnets are compensated by the U-arrangement of the stator, so that the guidance of the actuator over the runner or leaf spring bearing arrangement is guaranteed. The use of solely non-magnetic material prevents interference fields and disturbances from arising which can be traced back to the use of iron parts positioned at a distance. In addition, the motor, in particular the runner, becomes significantly lighter.\nFinally, it is advantageous that the U-formed linear motor is aligned in relation to a level which is fixed by the actuator parallel to a level which is formed by the floor space used by the base plate. The application of force by the actuator on the runner is thus free of moments, since the distance between the point at which the force is applied and the centre of gravity of the runner is almost zero in the vertical direction."}
{"text": "Ball bats, such as baseball and softball bats, are well known. In recent years, metallic bats including a tubular handle portion and a tubular barrel portion have emerged providing improved performance and improved durability over crack-prone wooden bats. The most common tubular bat is the aluminum single-wall tubular bat. Such bats have the advantage of a generally good impact response, meaning that the bat effectively transfers power to a batted ball.\nGenerally speaking, bat performance is a function of the weight of the bat, the size, and the impact response of the bat. The durability of a bat relates, at least in part, to its ability to resist denting and depends on the strength and stiffness of the tubular bat frame. While recent innovations in bat technology have increased performance and durability, most new bat designs typically improve performance or durability at the expense of the other because of competing design factors. For example, an attempt to increase the durability of the bat often produces an adverse effect on the bat's performance.\nThe incorporation of these advances and the use of additional materials, such as, other aluminum alloys, titanium alloys and composite materials have resulted in a large variety of well-performing ball bats. A typical metal bat, such as an aluminum bat, is formed with a one piece integral frame. Recently, high performance bats, such as bats incorporating the DeMarini® Half and Half™ bat technology, have been formed with separate handle and barrel portions, wherein the handle portion can be formed of a first material, such a composite material, and the barrel portion can be formed of a second material, such as a metal or a different composite material.\nOne drawback of recent ball bats formed of aluminum, titanium or composite materials is their cost. Aluminum, titanium and composite materials generally have a high material cost. For example, aluminum can cost up to ten times the price of conventional steel, and titanium is significantly more expensive than aluminum.\nDespite such advances in ball bat design and materials, a continuing need exists to further improve and optimize the performance, durability, feel and appearance of existing bats. It would be advantageous to optimize the weight distribution of a ball bat by removing or transferring material from one or more locations on the bat and redistributing, some or all of the weight of the removed material, to other more desirable locations. A need exists for design features that reduce the cost of a high performance ball bat without negatively affecting performance or durability of the ball bat. It would also be advantageous to produce a high performance ball bat with innovative design features that significantly improve the appearance of the ball bat without negatively affecting the performance of the ball bat. A need also exists for a ball bat that provides the batter with enhanced feedback during use."}
{"text": "1. Field of the Invention\nThe invention relates to a sliding window and, more particularly, to one which is adapted for use in automobiles and is raised and lowered electrically, and is equipped with a safety device and a height detector.\n2 Description of the Prior Art\nSliding windows used, illustratively in automobiles and equipped with safety devices, have been described in U.S. Pat. No. 3,651,389 and in German Pat. application No. P 29 39 942. As described therein, a capacitive pickup comprising a conductive strip is placed on the surface of the window and runs along its upper edge. Connected to this pickup is a conductive connecting line which is oriented along a vertical or oblique edge of the window to a point where a connection is made to a cable which connects the pickup to an electronic proximity detecting circuit.\nProximity safety devices of this sort should respond even if the pickup does not come in direct conductive contact with a part of a human's body. This occurs, for example, when the edge of the glass is touched with, for example, a gloved hand. In order for the safety device to respond under this condition, the proximity detector should have a very high sensitivity, i.e., a very high gain. However, if the detector is set to such a gain, adverse affects, such as circuit instabilities and spurious signals, which are unrelated to the approach of a part of the body, will cause the detector to produce an erroneous indication. In addition, windows having safety devices of this sort cannot be mass produced by known standard processes."}
{"text": "Illumination devices are useful in many areas of modern life, such as illumination of working area, illumination for the purpose of signaling and illumination carrying different type of information, such as an image and the like.\nIn medical applications, the requirement for good illumination while operating, both in the operating room (OR) and in emergency and pre-hospital situations is one of the most critical elements of successful surgery and treatment. Until now, lighting techniques in have evolved around the use of strong ceiling and head mounted light fixtures. Even while using minimal invasive techniques, illumination and good access for light source are a critical element is the treatment's success.\nUnfortunately, the physical bodies of the surgeons, doctors, paramedics, nurses and devices tend to block light when delivered from ceiling mounted light fixtures, both reducing the overall amount of light available, and producing confusing shadows. Head mounted light fixtures, on the other hand, may provide more directed light which avoids some of the problems of ceiling fixtures, but create other problems, such as strain on the head and neck of the surgeon or nurse wearing such a light.\nHand-held lights, typically held by a nurse, may also be used, but result in other problems, such as the requirement for a nurse or other medical personnel to hold and operate the light, and the necessarily limited size and weight of such lights. Also, this light is directed to the place the surgeon is looking at, not to the place where it is needed. Also, for all such lights, there is a problem of the difference in the quantity of light provided between an area that may be lighted by a hand-held or other light, and areas that do not receive light. This contrast or difference further increases the need for a supplemental light to reduce the differences.\nOf course, all of these different types of lighting are best suited for OR and/or emergency room use in controlled hospital settings. Yet even in such an environment, such lighting cannot always provide sufficient illumination for a cavity. This is because of obstacles such as the required angle of incidence and/or various structures that may block the light such as the organs or tissue of the patient as well as the surgeon's hands.\nEmergency and/or surgical medicine in some applications, such as military environments for example or pre-hospital situations (trauma center, field medicine, airborne evacuation) provides an even greater challenge as it must be performed under less than optimal medical conditions. These different types of lighting are not always suited for such conditions.\nOne attempt to overcome this problem has been to combine a cannula, which is a surgical tool, and a light-source. For example, Carlson et al. in U.S. Pat. No. 5,569,254 claim a surgical resection tool with irrigation, lighting, suction and vision. However, it is not disposable and light is conducted through a dedicated optical fiber.\nIn U.S. Pat. No. 5,425,730, Luloh describes a special cannula for eye surgery that has a plurality of optical fibers running outside said cannula externally.\nBerkowitz et al. in U.S. Pat. No. 4,551,129 claim a technique and apparatus for microsurgery including lighter-irrigator hypodermic tube. Light is conducted through optical fibers.\nIn U.S. Pat. No. 3,626,471 Florin designed a suction/washing unit for brain surgery that is comprised of two tubes having a fiber optic light-source at their front-end.\nA suction and illumination device is described in U.S. Pat. No. 3,261,356 by Wallace, in which light is provided through optical fibers located between two concentric tubes.\nSchultz in U.S. Pat. No. 5,281,134 provides light to various dental instruments by means of a single continuous optical fiber.\nU.S. Pat. No. 5,152,686 to Duggan and Jennings describes a dental appliance that includes a bite block and a suction tube attachment. A fiber optic light source is slidably and removably secured within the bite block.\nU.S. Pat. No. 5,139,420 to Walker claims a dental mirror system having a fluid conduit and/or a suction unit and a light transmitting cable to illuminate the work area.\nU.S. Pat. No. 3,995,934 to Nath describes a flexible light guide of the liquid filled type. The liquid is contained in a flexible plastic tube. An infrared light is guided through the liquid within a living body for medical applications.\nU.S. Pat. No. 5,580,154 to Coulter et al. discloses a glove apparatus including a glove member treated with an illuminative substance having phosphorescence or fluorescent illuminative properties, and a light circuit system integrally packaged therewith. Light from the light circuit system is contained within a ring-like light housing member which is preferably mounted on a finger section of the glove member. The outer glove surface may be decorated to ornately represent a fictional cartoon character, when the glove is implemented as a toy item.\nWO 97/31219 to Trow discloses a work or surgical glove and illuminator assembly which includes an illuminator oriented to project a light beam distally of the glove toward the work surface. The illuminator may have a self-contained light source, or utilize fiber optic-transmitted light from a remote light source. The illuminator may be disposed within the interior of the glove and projects a beam of light through a glove tip which is translucent or transparent. The assembly is useful when examining or operating upon an anatomical part of a patient.\nGB 2343361 to Spooner discloses a glove, particularly for cyclists, designed to give an illuminated signal so that other road users are made aware of the cyclist's presence. The glove includes an electric light, an electric battery and two electrical contacts arranged to be touched together to switch the light on by completing an electrical circuit between the battery and the electric light. In use, the wearer of the glove brings his thumb and forefinger together so that the two electrical contacts touch and the electric light is illuminated.\nU.S. Pat. No. 5,535,105 to Koenen relates to a glove which has a source of illumination mounted on the glove itself, for projecting light through a glove tip.\nGerman Patent No. DE 19952430 is also of interest as background art.\nIn all aforementioned patents, there is a source of illumination mounted on the glove or other device, but not incorporated within the glove. It should be noted that some of these patents indicate non-medical uses for the glove or other device, as such devices clearly have a broad range of possible applications, including but also extending beyond medical applications.\nIssalene and Lantrua in U.S. Pat. No. 4,872,837 claim a surgical or dental instrument and cannulae for aspirating, cleaning, drying and illuminating. A light source is placed in a housing behind a sleeve. The cannula is shaped and positioned for the light generated by the light source to be conveyed in the sleeve then in the cannula, whose wall performs the role of a light guide. Light is conducted axially to the cannula's open end where it emerges out. In one of the embodiments, the cannula tip has a circular bead that improves light distribution at the cannula's tip. However, the cannula's front tip, that is the only light emitting zone, provides a very limited illuminated area. Furthermore, the cannula's tip may be very close to or even in contact with soft tissues that may block the light.\nSuch “tip-only” illuminating cannulae provide a limited illumination area. In many medical procedures it is advantageous to have a much larger area being illuminated and seen clearly. Also, cannulae in general can only provide light in a limited area near the cannula itself, regardless of the overall area being illuminated. This is a clear disadvantage, because the functions of the cannula may require it to be located at a distance from the area which must be illuminated for the surgeon and/or other medical personnel. Also, a cannula may not always be present.\nFurthermore, such cannula are clearly only useful for surgical procedures. A more advantageous lighting device and system would provide directed light for many different types of procedures and environments, including non-medical environments."}
{"text": "1. Field of the Invention\nThe invention relates to digital data processing circuits. In particular, the invention relates to the performance of data manipulation functions performed on strings of data elements.\n2. Description of the Related Art\nConventional microprocessing circuits include several common building blocks. Essentially all such systems include a main memory storage area for storing data and instructions, and an execution unit for operating on the data in accordance with the instructions. After the function specified by a given instruction is performed, processed data is returned to the main memory storage area.\nIncreases in processor performance have been obtained by enhancements to this fundamental scheme. The processor may include two or more separate execution units which can process multiple instructions in parallel. The Intel Pentium and Pentium Pro are two examples of this type of processor. In some cases, different execution units are dedicated to different functions. The Intel Pentium Pro, for example, includes separate execution units for floating point and fixed point arithmetic operations. Another performance enhancement in almost universal use is the provision of data and instruction caches which provide local storage of recently used data and instructions. This speeds the fetching and storing of data and instructions by reducing the number of memory accesses required from a typically much slower main memory storage area.\nStill, some types of operations are performed inefficiently by these processor architectures. One inefficiently performed class of instructions is data string manipulation instructions. In these instructions, operations involving a sequence of data elements are performed. For instance, a block of data may be moved from one series of memory addresses to another series of memory addresses. Alternatively, the elements of a block of data may be compared to a test data element or a string of test data elements. In the Intel Pentium Pro, assembly language instructions are provided to perform these functions on a specified string of data. Although the total length of the processed string can be very large, data is moved and/or analyzed as short string portions of at most 32 bits long due to the bus width and 32 bit execution unit. Performing a string move on the Pentium Pro thus involves a sequential process of reading and writing pieces of the data string to and from main memory (or to and from the cache for those portions of the string which are present there). String scans for matching data are similarly performed. Short pieces of the string are read from memory and compared to the desired test string. Thus, the Pentium Pro architecture includes useful string instructions as part of its instruction set, but is incapable of performing string operations on large strings as quickly as memory technology might allow.\nAnother commercially available device which includes string manipulation features is the TMS320C80 digital signal processor from Texas Instruments. Because this device is adapted for use in video and multimedia applications, features for improving the speed of movement of large blocks of data such as a set of image pixels are included. In the TMS320C80, the programmer may write string movement parameters to a memory location. These string movement parameters can then be transferred to the memory controller portion of the device, and the string movement is performed by the memory controller without further involvement of the execution unit. This feature helps speed up the movement of data blocks, but setting up the transfer parameters requires preliminary write operations, which is inconvenient for the programmer, and results in the need to use several instructions to initiate a block move. Furthermore, although the TMS320C80 includes a data cache, these memory move operations are not performed utilizing cached data, and no mechanism to ensure cache coherency with the main memory where the data move occurs is provided.\nOther implementations of memory systems which can perform data manipulation have been described. In U.S. Pat. No. 5,590,370, a system is disclosed which includes “active memory elements” that incorporate processing logic for performing searches and other data manipulations outside of the host processing circuit. U.S. Pat. No. 4,731,737 also discusses memory elements which can receive data manipulation commands from an external host processor. However, neither of these systems provide for cache coherency, and they further do not describe the use of assembly language instruction sets which provide simple and efficient programming of data string manipulations. Thus, there is a continuing need for improvements in processor architectures to allow the processing of data strings quickly and efficiently."}
{"text": "In internal combustion engines, fuel is transported by means of a fuel pump from a tank in preparation for a subsequent high-pressure pump. The high-pressure pump is normally driven by the internal combustion engine and transports the fuel into a fuel pressure accumulator (fuel rail). The high-pressure pump itself is not regulated and conveys the fuel which is made available at its inlet connection into the fuel pressure accumulator.\nIn order to provide a defined fuel quantity to the high-pressure pump, provision is made for a volume flow control valve between the fuel pump and the high-pressure pump, said volume flow control valve being controlled by a control unit. The fuel through-flow through the volume flow control valve is adjusted depending on a current which flows in a valve coil of the volume flow control valve. The pressure in the fuel pressure accumulator can be adjusted via the fuel volume which is provided to the high-pressure pump.\nThe volume flow control valve normally has a leakage flow in the zero-flow state. This can result in an unwanted fuel pressure increase in the fuel pressure accumulator when the injection volumes are very small or, for example, no fuel is injected in the case of an overrun cut-off.\nAs a result of the construction, preventing the leakage flow in the zero-flow state of the volume flow control valve can only be achieved at considerable cost and is moreover undesirable in specific cases, if emergency operation of the internal combustion engine occurs in the event of a failure of the volume flow control valve or the control unit.\nA regulator valve is normally provided at the fuel pressure accumulator, by means of which regulator valve the pressure in the fuel pressure accumulator can be adjusted depending on a control current. The regulator valve is actively controlled by the control current, such that the pressure in the fuel pressure accumulator is adjusted depending on the control current and depending on a fuel flow through the regulator valve. The fuel flow must exceed a threshold value so that the regulator valve can be operated in a linear range. This additional fuel flow through the regulator valve must be conveyed by the high-pressure pump, in order that the regulator valve can be operated in the linear range. When the high-pressure pump is dimensioned, it is therefore necessary to ensure that the high-pressure pump supplies the regulator valve with a minimum through-flow and moreover provides the volume that is required to build up the pressure or to maintain the fuel pressure in the fuel pressure accumulator."}
{"text": "The present invention relates to a disk player which regenerates information from a disk-type recording medium and/or records information on the same, and more particularly to a disk player of a kind, which is improved in relationship between opening/closing timing of a shutter arranged on a disk cartridge for accommodating the disktype recording medium, and an information regenerating location and/or an information recording location. The present invention also relates to a disk player of a kind, which is improved in a driving mechanism thereof, for upward and downward shifting a recording mechanism which records information on the disk-type recording medium.\nAs is widely known, a disk cartridge of a disk player, which accommodates therein a disk-type recording medium, is provided with a slidable shutter for properly exposing an information surface of the disk-type recording medium. The shutter is closed normally, to thereby protect the disk-type recording medium in the disk cartridge from an external environment.\nTherefore, in conventional disk players, when information is regenerated from the information surface of the disk-type recording medium or the information is recorded on the same, the following complicated operations are required: First, the disk cartridge is carried to a predetermined location by a carrier mechanism, and at the same time the shutter of the disk cartridge is fully opened by means of a shutter opening/closing mechanism mounted in the disk player, to thereby expose the information surface of the disk-type recording medium. Then, the information is regenerated from the exposed information surface by means of an optical pickup provided in the disk player, or the information is recorded on-,the exposed information surface by means of the optical pickup and a magnetic head.\nThe above-mentioned conventional disk players are operated in the following two modes with respect to information regeneration and information recording:\nIn the first mode, when the disk cartridge is located at an information regenerating location in the disk player, the shutter is fully opened by means of the shutter opening/closing mechanism. Thereafter, while the disk cartridge is further shifted from the location to an information recording location in the disk player, the magnetic head is disposed in sliding contact with the information surface of the disk-type recording medium, to thereby carry out information recording.\nIn the first mode, however, when the disk cartridge is further shifted from the information regenerating location to the information recording location, an additional configuration for shifting the shutter opening/closing mechanism of the disk player together with the disk cartridge, or alternatively a configuration for detaching the shutter opening/closing mechanism from the disk cartridge in the above state while maintaining the fully open state of the shutter is required. Therefore, the configuration for shifting both the carrier mechanism and the shutter opening/closing mechanism or the configuration for separating the two mechanisms from each other complicates and upsizes the disk player, resulting in an increased number of component elements and hence a rise in manufacturing costs.\non the other hand, in the second mode of the disk player, the information regenerating location and the information recording location are set to the same location, and the shutter is fully opened at the location, followed by disposing the magnetic head in sliding contact with the information surface of the disk-type recording medium to carry out information recording.\nIn this mode, however, the carrier mechanism must be detached from the disk cartridge, and then the magnetic head is driven to be disposed in sliding contact with the information surface of the disk-type information recording medium, by the carrier mechanism. Alternatively, a magnetic head driving mechanism for driving the magnetic head must be additionally installed. Therefore, if the carrier mechanism is further shifted or the magnetic head driving mechanism for driving the magnetic head is additionally employed, the construction of the disk player is further complicated and upsized. In addition, the number of component elements is increased, which inevitably leads to increase manufacturing costs of the disk player.\nBesides, in either mode, to hold the information recording medium inserted into the disk player, on a turn table, a supporting member for supporting the turn table need to be relatively shifted with respect to the information recording medium, and further, during information recording, the magnetic head need to be relatively shifted with respect to the information recording medium. Thus, both the supporting member and the magnetic head need to be relatively shifted with respect to the information recording medium, which complicates and upsizes the driving mechanism, resulting in a further inconvenience, e.g. that the accuracy of a vertical location cannot be improved."}
{"text": "The batteries (main batteries) in a passenger aircraft are important components of the on-board power supply system. For decades, rechargeable nickel-cadmium batteries having a nominal voltage of 24 volts have been used for this purpose, which nickel-cadmium batteries consist of twenty nickel-cadmium battery cells which are connected in series and each have a nominal voltage of 1.2 volts. The batteries are used for example to supply power in emergencies until other redundant power sources become available, to start the auxiliary power unit (APU), and to supply power to sub-systems on the ground for servicing and in preparation for flight. In these cases, strict requirements are placed on safety, reliability and availability, as a result of which servicing intervals are relatively short.\nAt present, the batteries are removed after a certain number of flight hours (FH) or after a defined period of time has elapsed. The batteries have to be removed from the aircraft and reconditioned in a workshop in order to determine the condition thereof and to carry out maintenance thereon.\nHowever, in their documentation, battery manufacturers always give operational hours (OH) for the servicing intervals. However, since aircraft operators only record flight hours, a factor is defined using the following ratio: factor=OH/FH. This ratio can vary within a particular range depending on the mission profile of the aircraft. The mission profile is calculated using the FH/FC ratio, where FC is the flight cycles, and is classified as being short haul, medium haul or long haul. A relationship to the turnaround time is thus established.\nA typical cycle of an aircraft battery, which is removed after the applicable time limit has elapsed, is for example storage-installation-operation in aircraft-removal-maintenance-storage. Measures taken during the maintenance stage are for example: periodic checks, including for example insulation measurements, voltage measurements and filling level measurements; regular checks including periodic checks and capacity tests; and general overhaul including regular checks and cleaning.\nIn an aircraft, it is conventionally only the total voltage and the charging or discharging current that is recorded. Furthermore, in some types of batteries, a temperature is also measured at a central measuring point. These data are not analysed or combined in the aircraft, it being impossible in any case to reach a satisfactory conclusion on the state of ageing of the battery on the basis of these data.\nIn the workshop, the state of health (SoH) of an aircraft battery, as set out by the manufacturer in the component maintenance manual, is determined using the insulation values, the cell voltage during loading, the filling level of the electrolyte and the capacity. There are doubts as to the validity of this method in respect of the state of health. If the aircraft battery passes the capacity test without any difficulty, it is deemed to be in working order irrespective of the operating condition of individual cells and can thus be installed in the aircraft. In this way, weakened cells go undetected and sustain further damage during continued operation. If the fault is later detected during maintenance operations, the battery should be disposed of if there is a low number of faulty battery cells, for example five. In any case, weak points or potential faults in a battery which is in use are not detected or remedied until maintenance is carried out in the workshop, if at all. These faults are thus dormant faults; it is not possible to detect these faults at an early stage or in a timely manner. Therefore, no preventative measures can be taken either in order to counteract irreparable damage and the associated requirement of permanently taking individual battery cells or the whole battery out of operation.\nIn the worst case scenario, the capacity required in an emergency, which should ensure that there is an emergency power supply for 30 minutes, cannot be drawn from a battery that has been previously damaged in this manner, and this can have an adverse effect on flight safety.\nHistorical analysis demonstrates that nickel-cadmium batteries used in aircraft are very reliable components. The number of times that the batteries are removed when undamaged is therefore relatively high, resulting in unnecessary effort. An analysis of a particular type of passenger aircraft carried out in 2012 showed for example that only a particular proportion of the removed batteries had a dormant fault which had not been detected during operation in the aircraft. Therefore, a fault was not found in many of the removed batteries, and therefore said batteries would not have needed to be removed.\nIf a battery cell weakens steadily over time, the other battery cells are loaded more heavily and may likewise become damaged. This imbalance between the individual battery cells cannot be detected from the outside. Even the tests required by the manufacturer only inadequately identify weakened battery cells. This increases the need to take whole batteries out of operation unnecessarily because a plurality of battery cells may already have a serious defect at the time when the fault is discovered.\nAnother drawback is that there is no universal standard for the removal of batteries. Each operator establishes, in consultation with system developers, the servicing intervals after which a battery should be removed. Empirical values and manufacturer recommendations form the basis of this decision. A solid database in the form of continuously collected measurement values does not currently exist. This poses a risk in respect of the evaluation of the state of ageing of the batteries.\nThe object of the invention is that of providing a system and a method for monitoring an aircraft battery which provide for cost-effective and efficient servicing and maintenance, increase the average operational life of the batteries and flight safety, and reduce the amount of irreparable damage to the batteries.\nThis object is achieved according to the invention by the features of the independent claims."}
{"text": "The delivery of digital television to the home was launched in earnest in 1995 for both cable television as well as satellite-delivery systems. This new technology enabled multi-channel video program distributors (MVPD) to provide far more television programming using available network bandwidth as compared to what had been possible using analog signals of the same programs. A plurality of digital television signals are combined to fit multiple digital channels into the space of one legacy analog channel via a process called “multiplexing.” When television programs were digitally encoded at a fixed bit rate (called Constant Bit Rate or ‘CBR’), then the early digital cable TV systems could carry perhaps six to eight digital television programs in the space of a single legacy analog channel (6 MHz for NTSC or 8 MHz for non-NTSC-based systems).\nThe distribution networks of MVPD systems, whether cable TV or satellite, are known as “managed networks” because the output of a multiplexer is typically of a fixed bit rate. For comparison, the Internet data network is known as an “unmanaged” network, since the public use of the Internet is not regulated by a central controlling mechanism and bandwidth between two points on the network varies unpredictably.\nVariable bit rate (VBR) video encoding is more efficient in the use of bandwidth than CBR encoding. VBR also generally delivers a better quality picture for the same average bandwidth. However, VBR is more difficult to manage on a distribution network. Statistical multiplexing is used to address this difficulty.\nWith the advent of interactive services hosted in a central location, such as a cable TV headend as well as with media originating “in the cloud” and routing over a managed network to a consumer set-top box, the difficulty of managing the VBR session within a multiplexer becomes far more challenging and more prone to adverse interactions among the sessions within a multiplex stream.\nInteractive television services provide the viewer with the ability to interact with their television for the purposes of selecting certain television programming, requesting more information about the programming, or responding to offers, among many possible uses. Such services have been used, for example, to provide navigable menu and ordering systems that are used to implement electronic program guides and on-demand and pay-per-view program reservations without the need to call a service provider. These services typically employ an application that is executed on a server located remotely from the viewer. Such servers may be, for example, located at a cable television headend. The output of a software application running on such servers is streamed to the viewer, typically in the form of an audio-visual MPEG Transport Stream. This enables the stream to be displayed on (or using) virtually any client device that has MPEG decoding capabilities, including a “smart” television, television set-top box, game console, and various network-connected consumer electronics devices and mobile devices. The client device enables the user to interact with the remote application by capturing keystrokes and passing these to the software application over a network connection.\nAn interactive television service combines the properties of managed and unmanaged network topologies. Such services require low delay, perceptually real-time properties typically associated with Real Time Transport Protocol running over User Datagram Protocol (UDP/RTP) high-complexity, proprietary clients. However, in interactive television applications the stream must be received by relatively low-complexity clients using consumer electronics-grade components. Typically, these clients do not have the capability of much more powerful laptop and tablet computers to which the user has grown accustom. Hence, interactive applications hosted on a cable or satellite set-top box are perceived as slow and old-fashioned compared to the contemporary norm. Hosting the application in a central means (e.g., a remote server located at a cable television headend) and providing the picture output to the set-top device mitigates this shortcoming and allow for the delivery of rich, highly interactive applications and services. It also places stronger demands on the distribution network to deliver these services.\nA centrally (remotely) hosted interactive television service provides a combination of relatively static image portions representing a Graphical User Interface (graphical UI or GUI) that requires low-latency, artifact-free updates responsive to user input, and other portions that may have video with associated audio that require smooth and uninterrupted play-out. Conventional network distribution systems do not adequately facilitate this combination of data types. For instance, with existing statistical multiplexers for cable or satellite television systems, when large user interface graphics of a particular session need to be sent to a particular client, the many other sessions sharing the same multiplex have no means available (except a drastic reduction in image quality) to scale back the bandwidth requirements of adjacent streams to allow a temporary large data block representing the UI graphics to pass.\nWith many interactive sessions active within a single multiplex stream, a possibility exists for disruption to video, audio and/or GUI data. The only alternative that conventional systems have is for the conservative allocation of bandwidth which then supports many fewer simultaneous sessions per multiplex stream.\nTherefore, it is desirable to provide an improved way for multiplexing interactive program streams.\nAdditional background information is provided in U.S. patent application Ser. Nos. 12/443,571; 13/438,617; and 14/217,108, all of which are incorporated by reference herein in their entirety."}
{"text": "A well-known embodiment of a stack arrangement provides for stacking two memory chips within an FBGA package. Here, the lower memory chip, with its active surface area face-down, is integrated into the circuit by means of the conventional Board-on-Chip (BoC) technology and the upper memory chip, with its active surface area face-up, is attached to the lower memory chip and integrated by means of wire bond connections.\nSince the electrical connections of the upper memory chip on the carrier substrate must be on the same side as those of the lower memory chip, this stack arrangement for memory chips with bond pads in a double-row arrangement, because of the simultaneous face-up and face-down arrangement in one stack, requires that the individual wire bond connections for the upper memory chip be crossed. In addition to prolonging the signal transit times, this design leads to technical problems because of the potential mechanical, electrical or thermal interactions between the wire bond connections at the crossover points. Another serious disadvantage of this stack arrangement is that bum-in and testing are possible only once the stack arrangement has been completely assembled, which leads to a relatively high loss rate which amounts to approximately the square of the loss rate of one discrete component.\nThus, one problem that can be solved by some embodiments of the present invention is to make available a stack arrangement for a memory module that meets the requirements of small dimensions and low electrical parasitism, approximately identical to that of discrete components, makes it possible to achieve especially low signal transit times, reduces the loss rate of the assembled arrangements, and can be integrated into the existing manufacturing and testing technology."}
{"text": "Reprogrammable logic devices, such as FPGAs (field programmable gate arrays), are devices that are configurable to perform a certain logic function. Such devices generally comprise millions of static random access memory (SRAM) cells, each storing configuration data for configuring the logic function to be implemented.\nRadiation, for example in the form of cosmic particles, can create errors in electronic devices, and are particularly problematic in the case of reprogrammable logic devices. Indeed the radiation may cause errors not only in the data that is being processed, but also in the configuration data, resulting in a change of the logic function implemented by the circuit.\nUS patent publication U.S. Pat. No. 7,764,081 proposes a solution for protecting a reprogrammable logic device from errors due to single event upsets (SEUs). According to this solution, the SRAM cells that store the configuration data in the reprogrammable logic device are replaced by DRAM (dynamic random access memory) cells, which are periodically refreshed from a PROM (programmable read-only memory).\nHowever, a problem with the solution proposed by U.S. Pat. No. 7,764,081 is that it is relatively complex and energy consuming. There is thus a need in the art for a simple and low energy consuming solution for protecting a reprogrammable logic device from errors caused by radiation."}
{"text": "In general, magnetic disk drives utilize read and write channels to write and retrieve information from magnetic medium. In order to improve the recording density of the data storage system enhanced equalization and data detection systems are sometimes employed. One such scheme is the use of partial response equalization and Viterbi decoding or detecting in the read channel.\nThe read channels for magnetic disk drives sometimes employ a strategy of peak position detection with simultaneous amplitude qualification. Such read channels for magnetic disk drives are suitable if the inter symbol interference (ISI) is largely removed by an equalization system and if the signal to noise ratio (SNR) is adequate following such equalization. The adequacy is a function of the tolerable error rate within the overall read system. For example, if error correction capability increases, then the SNR may be decreased while still providing adequate overall system performance. In such reading systems, noise is boosted by equalization thereby decreasing SNR. As recording densities increase, the inter symbol interference becomes larger and it contains a larger portion of non-linear ISI which cannot be effectively removed by equalization. In any case, SNR decreases causing an overall degradation in error rate performance.\nSome magnetic disk drives employ a technique known as partial response maximum likelihood (PRML) detection. In most cases, the PRML detection systems utilize some form of Viterbi detector. In general, density gain is achieved in PRML in two ways: Viterbi detectors tolerate larger amounts of ISI thereby minimizing equalization noise boost, and Viterbi detectors are inherently more robust than peak detectors.\nOne problem associated with utilizing PRML systems is achieving well controlled partial response pulse shapes at the output of the equalizer. The ability to achieve well controlled pulse shapes at the output of the equalizer is typically dependent upon the existence of mostly linear ISI (i.e. created from so called linear superposition of pulse shapes), and it requires sufficiently complex equalizers that generate the partial response pulse shapes. The classical Viterbi detector depends upon discrimination between amplitude levels, which are characteristic of linear ISI, from superimposed partial response pulse shapes.\nThe use of high linear densities in magnetic recording systems results in the ISI having a large non-linear component. Consequently, classical PRML error rate performance is thereby significantly degraded below the level of performance obtained if the ISI was entirely linear. Furthermore, the cost in logic complexity and power consumption for constructing long finite impulse response (FIR) equalizers to create the best possible pulse shapes is an undesirable burden. Therefore, compromises are typically made in equalizer complexity, including both length and precision, to limit the amount of logic and power consumption. It is highly desirable to develop an equalizer and detector system that is less demanding in terms of perfection in equalizer output.\nThe PRML systems have been used in magnetic disk drives as well as other applications. One such system for disk drives uses a PR IV (partial response, class IV) pulse shape and a 0,4 recording code. In such PRML systems, the Viterbi decoder was matched to the particular pulse shape and code selected. The PR IV and 0,4 recording code permitted a simple Viterbi decoder with excellent bandwidth. However, no compensation for non-linear ISI was incorporated.\nOther disk drive systems have employed more elaborate PRML systems based on a 1,7 recording code and selectably either an EPR IV (extended partial response, class IV) or E.sup.2 PR IV pulse shapes. These disk drive systems included a Viterbi decoder to match the equalizer. The equalizer contains taps added to partially correct for pole tip undershoots, which are characteristic of widely used thin film recording heads. However, the system does not compensate for non-linear ISI except by use of a modest degree of target shape compensation."}
{"text": "As used herein, the terms “user equipment” and “UE” in some cases might refer to mobile devices such as mobile telephones, personal digital assistants, handheld or laptop computers, and similar devices that have telecommunications capabilities. A UE might include a device and its associated removable memory module, such as but not limited to a Universal Integrated Circuit Card (hereinafter, UICC) that includes a Subscriber Identity Module (hereinafter, SIM) application, a Universal Subscriber Identity Module (hereinafter, USIM) application, or a Removable User Identity Module (hereinafter, R-UIM) application. Alternately, such a UE might include the device itself without such a module. In other cases, the term “UE” might refer to devices that have similar capabilities but that are not transportable, such as desktop computers, set-top boxes, or network appliances. The term “UE” could also refer to any hardware or software component that could terminate a communication session for a user.\nMobile communications has become an everyday commodity. In a communication system, radio resources are shared by a plurality of control channels and a number of data channels between base stations and user equipment. A base station identifies UEs through control channels and transmits data to the UEs via data channels. When the number of UEs connected to the base stations is relatively low, system resources are sufficient for control channels and data channels of all connected UEs. However, number of UEs connected to the base station varies a lot in different environments.\nFIG. 1 is a schematic diagram illustrating a scenario of machine type communications. The numerous UEs 11 access and share the system resources simultaneously. For the machine type communication, the demands for higher control channel capacity and data channel capacity dramatically increase. Therefore, how to multiplex the control channels and the data channels corresponding to UEs 11 wherein the improvement of flexibility and efficiency of the resource utilization is a worth research issue."}
{"text": "The present invention relates to electronic devices and methods for specifying a time when a specified process was performed, and transmitters and receivers for use in the electronic devices.\nConventionally, some electronic devices such as facsimile devices/pagers are known which record/store the time when they receive/transmit data from/to other electronic devices. In the case of the electronic device which stores the time when it transmits data, the time output from a clock circuit provided in the electronic device thereof is used whereas in the case of the electronic device which stores the time when it receives data, the time output from a clock circuit provided in the electronic device is used.\nThere are many demands for transmission of some obtained target data along with data on the time when the target data was received (and not along with data on the time when the target data was transmitted), when the time data is significant. To this end, the obtained target data and the data on the time when the target data was obtained are temporarily stored in a memory and then transmitted to the requesting end.\nFor example, in the case of an in-home examination and treatment where a patient measures his or her living body data such as his or her pulse, blood pressure or electrocardiogram waveform at his or her home, and sends the data over a telephone line to medical facilities, it is necessary to further send data on the time when the living body data were measured. If otherwise, it cannot be known when the target data were obtained, and the target data cannot serve as a good reference for diagnosis.\nIn addition, for example, a time when an image was picked up by an electronic still camera having a communication function can be desired to be sent along with the picked-up image.\nWhen data on the time when such data to be sent is stored and sent, a clock circuit is required to be provided at the transmitting end. In this case, the time recorded by the clock circuit is required to be accurate at all times.\nIf a clock circuit is also provided at the receiving end, the sending and receiving ends are required to be synchronized. Thus, in order to satisfy such requirement at all times, the transmitting and receiving ends are continuously required to be synchronized or corrected with respect to time, which is very troublesome.\nIf living body data measuring devices are required for the respective measuring items, for example, if separate devices such as electrocardiograms and pulse meters measure the corresponding data and send same, in an in-home examination and treatment, the respective devices are required to be synchronized, which is very troublesome.\nIt is therefore an object of the present invention to provide electronic devices which include a transmitter and a receiver, and which are capable of storing and specifying a time when the data to be transmitted was obtained even when no current time measuring means is provided in the transmitter.\nAnother object of the present invention is to provide electronic devices and methods which eliminate troublesome synchronization between the transmitter and receiver even when a current time measuring means is provided in the transmitter, and to further provide transmitters, receivers and storage mediums which are operable with such electronic devices."}
{"text": "The present invention relates to a novel and useful tackifier.\nBinders or tackifiers are employed to bind together many kinds of materials, especially particulate matter. For example, tackifiers may be employed as dust palliatives, landfill stabilizers, alternate daily cover, binders for manufacturing and mining waste products, a component for hydraulic planting, straw tacking, erosion controllers, soil stabilizers, components in fire retardant and insulative materials, and the like. Hydro-mulching or hydro-seeding, straw tacking, alternate daily cover, as well as insulation applications, often require the mixing of fiber with a tackifier. Of course, hydro-seeding also includes the addition of particular seeds which are intended to grow from the soil being spread with the tackifier, fiber, fertilizer and seed combination. Tackifiers may be used alone for the above purposes.\nMany methods and products have been proposed for combining with soils or other particulates. For example, U.S. Pat. No. 4,865,640 describes a moisturizing substrate utilizing cellulosic compounds combined with hydrated metallic salts and aerated water through an ether linkage.\nU.S. Pat. Nos. 5,556,033, 5,849,364, and 6,096,373 describe sprayable compounds which are formed into a slurry and used to cover items forming a water repellant layer.\nU.S. Pat. No. 4,369,597 describes a process for making pelletized mint mulch derived from the mint plant. The mint pellets may be spread or sprayed by hydro-mulching machines, and may also contain a natural tackifier to adhere fibers to the soil surface.\nU.S. Pat. No. 3,763,072 shows an erosion control crust formed of acrylic latex emulsion and sodium silicate.\nU.S. Pat. Nos. 5,082,500 and 5,516,830 illustrate a sprayable composition of a polymer mixed with clay and water to provide a flexible water repellent layer over soil, refuse, or grain.\nU.S. Pat. No. 6,029,395 teaches a biodegradable mulch mat utilizing hair fibers, granulated corn, and a vinyl polymer as a binder. The mulch mat is used to improve soil temperature, suppress weeds, and build soil, while conserving soil and water.\nU.S. Pat. No. 4,705,816 describes a liquid mulch prepared from a polymeric binder and fine particulate filler particles.\nU.S. Pat. No. 6,076,299 delineates mulching pellets made from finely divided paper and wood, a water absorbent natural polymer, and a surfactant. Natural polymer may be a guar gum.\nU.S. Pat. No. 3,938,279 puts forth a growth medium which is used to cover a soil surface consisting of a mixture of cement, fiber, and fertilizer.\nU.S. Pat. No. 3,600,852 shows a seed containing composition utilizing gypsum and enough water to produce a slurry.\nU.S. Pat. No. 5,459,181 teaches a hydraulic binder composition having a suitable hydroxyl bearing polymer, a covalent cross-linking agent, and a catalyst. The final formulation is an insoluble structure and serves as a binder or tackifier.\nU.S. Pat. No. 5,801,116 describes a process for producing polysaccharides in a coarse particle size. The composition is used to absorb materials.\nU.S. Pat. No. 3,875,697 shows the cross-linking of acrylamide polymers with hypochlorite and hypobromite salts of alkali and alkaline-earth metals.\nA binder and tackifier which is easily prepared and delivered to a site through spraying techniques known in the art would be a notable advance in the field of particulate control.\nIn accordance with the present invention a novel and useful tackifier is herein described.\nThe tackifier of the present invention utilizes a synthetic polymer consisting essentially of polyacrylamide or polyacrylate, or a co-polymer of acrylate and acrylamide. Such polymers are commonly referred to as PAM. It is known that PAM may be used as a tackifier in 100% concentrations combined with water. However, PAM is quite expensive to employ in this regard. Also, it is known to spray solutions of natural organic materials alone, such as cornstarch for use as a binder or tackifier. However, such natural materials are highly susceptible to bacterial degeneration and are not suitable for long-term use. In addition, natural organic materials require application in high volumes, typically 60-120 pounds per acre.\nCombined with the synthetic PAM polymer of the present invention is a natural organic material. Such natural organic materials may take the form of cornstarch, rice starch, wheat starch, potato starch, polysaccharides, modified polysaccharides, mint, animal gelatin, plant seed husks, corn gluten, natural latex, guar gum, psyllium, mucilage, or kelp. Although the simple mixture of PAM and a natural organic material produces an association, theorized as a molecular entanglement, a cross-linking agent is also employed in the present invention to cross-link the PAM polymer to the natural organic material. For example, the cross-linking agent may take the form of sodium borate, sodium metaborate, sodium tetraborate, hydrated sodium borate, hydrated sodium metaborate, hydrated sodium tetraborate (borax), magnesium sulfate, molybdenum sulfate, or sodium molybdate.\nThe PAM, natural organic material, and cross-linking agent may be mixed together with a sufficient amount of water to form a sprayable slurry. Fiber, such as paper or wood fiber, may be added to the slurry to form an insulative or erosion control material. Fertilizer may also be added to the slurry to enhance seedling development. In addition, seed may be added to the fiber to form a fiber slurry which may be used for hydro-seeding purposes.\nIn certain cases, the tackifier combination of the present invention may include a surfactant to improve permeability and wetting of the tackifier of the present invention. Moreover, a colorant may be added as an indicator to show the extent of spreading of the tackifier of the present invention over a surface. In addition, admixtures may be formed with the tackifier of the present invention and water absorbent polymers to increase water-holding capacity of the same.\nIt may be apparent that a novel and useful tackifier has been described.\nIt is therefore an object of the present invention to provide a tackifier which may be used to control soil erosion, dust, and migration of particulate matter of various types.\nAnother object of the present invention is to provide a tackifier utilizing the relatively expensive synthetic organic component and a relatively inexpensive natural organic material to form an effective low cost tackifier which may be used for various purposes.\nAnother object of the present invention is to provide a tackifier using multiple synthetic and natural components which may be cross-linked into a macromolecule to add stability to the tackifier.\nAnother object of the present invention is to provide a tackifier which may be employed to control the migration of particulate matter, to aid in hydro-seeding, and form other layers which exhibit insulation and fire retardancy characteristics.\nAnother object of the present invention is to provide a tackifier including components that possess low toxicity, permitting safe handling, and are compatible with sources of water, including surface, ground, and potable water sources.\nA further object of the present invention is to provide a tackifier for various uses which is inexpensive to manufacture and is easily spread by conventional methods.\nAnother object of the present invention is to provide a tackifier for various uses which sets up in a relatively short time period, is relatively durable, and re-solubilizes after an initial drying.\nA further object of the present invention is to provide a tackifier which may be employed in hydro-seeding, hydro mulching, or hydraulic planting processes.\nThe invention possesses other objects and advantages especially as concerns particular characteristics and features thereof which will become apparent as the specification continues.\nVarious aspects of the present invention will evolve from the following detailed description of the preferred embodiments thereof which should be taken in conjunction with the appended examples.\nThe tackifier of the present invention utilizes a synthetic polymer known as polyacrylamide, polyacrylate, or a co-polymer of acrylamide and acrylate (PAM). PAM is produced in various molecular weights and is generally a synthetic material. The PAM component in the present invention varies between 10 and 80 percent by weight in an operating range. In certain instances the polymer may be present in a preferred range of 10 to 60 percent by weight. Finally, an optimum range for the synthetic polymer would lie between 50 to 65 percent by weight. The preferred form of PAM contains 0.05% acrylamide monomer or less to minimize toxicity.\nAnother component of the tackifier of the present invention would take the form of a natural organic material. Such natural organic materials may be considered to be waste materials and are relatively cheap to obtain. However, it has been found in the present invention that such natural organic materials in the presence of PAM forms a relatively durable tackifier either by co-polymerization and/or molecular entanglement. Such natural organic materials may include cornstarch, rice starch, wheat starch, potato starch, polysaccharides, modified polysaccharides, animal gelatin, plant seed husks, mint, corn gluten, natural latex, guar gum, psyllium, mucilage, and kelp. In operative combination with PAM, the natural organic material may range between 15 and 75 percent by weight. Likewise, the preferred range of the natural organic material lies between 15 and 45 percent by weight. Also, the optimal range for the natural organic material exists between 25 and 35 percent by weight.\nA cross-linking agent is also employed in the present invention. The cross-linking agent takes the form of an inorganic salt. It has been found that sodium borate, sodium metaborate, sodium tetraborate, hydrated sodium borate, hydrated sodium metaborate, hydrated sodium tetraborate (Borax), magnesium sulfate, molybdenum sulfate, and sodium molybdate, are sufficient in this regard. Other inorganic salts would also serve as cross-linking agents between the PAM and the natural organic material. The operative range of the cross-linking agent is 1 to 15 percent by weight. A preferred range of the cross-linking agent lies between 1 and 10 percent by weight. Also, the optimum range of the cross-linking agent ranges between 8 and 12 percent by weight.\nIn certain instances, a surfactant may be employed to aid in the tackifier\"\"s penetration and wetting characteristics. Such surfactant may be employed in an amount ranging between a trace amount and to about 1 percent by weight.\nTo permit the tackifier of the present invention to be sprayed or delivered through a nozzle, a sufficient amount of water may be added to the tackifier formulations delineated above. The latter addition produces a sprayable composition.\nIt should also be noted that the tackifier above described may be mixed with fiber of various kinds to produce insulative coatings, fire retardant coatings, and the like. In addition, seeds may be added, with or without the fibers, to produce a mixture which is suitable for hydro-seeding purposes.\nAlso, a small quantity of inert coloring material, less than 1 percent by weight, may be added to the add mixture during production and packaging.\nIt has been found that the tackifier of the present invention uses ⅓ less PAM than a tackifier employing pure PAM. This results in the tackifier of the present invention being useable as an effective, low cost, product. Moreover, the tackifier of the present invention may be employed in lesser amounts than natural organic materials. For example, 60-120 pounds per acre of natural organic material is required, versus about 3 pounds per acre of the tackifier of the present invention.\nThe following examples are submitted to illustrate the invention in further detail, but are not deemed to limit the invention in any manner."}
{"text": "1. Field of the Invention\nThe present invention relates to a power and signal transmission system which makes signal transmission via lines and permits transmission of power by using lines such as those used in a public telephone network.\n2. Description of the Related Art\nIn general, in exchanges, etc., power transmission is made via signal-transmitting lines. In a typical exchange, first and second communication devices are connected to each other via communication lines. A digital signal is applied to the input terminal of the transmitting circuit of the first communication device and is converted by a transmitting circuit, which is powered from a power source, to a signal having no direct current component. The converted signal is, in turn, transmitted on the lines via a signal-source resistor and a transformer. The receiving circuit of the second communication device receives the signal transmitted over the lines. In the first communication device, a high-voltage source delivers a high-voltage from the center tap of the secondary winding of a transformer, connected by a signal-short-circuiting capacitor, to the lines. The high voltage supplied to the lines is lowered by line resistance and then applied to a DC-to-DC conversion circuit in the second communication device. The conversion circuit converts the input high voltage to a power supply voltage required with a receiving circuit in the second communication device.\nHowever, the use of transformers and choke coils for power transmission will increase the cost and make the system large because they are costly and large. In addition, there are disadvantages in that the assembly of the transformers and the choke coils takes a lot of time and labor, and the assembly cost also increases."}
{"text": "Hematophagy is a behavior of feeding on blood taken from vertebrate hosts. It is practiced by about 15,000 species and about 400 genera of arthropods (Mans, Louw et al. 2002). It has also been reported that blood-feeding has evolved independently at least six times. In addition to arthropods, members of other phyla including annelids such as worms and leeches are hematophagous. While independent evolution and wide-spread occurrence of taking blood-meals have resulted in considerable biological diversity, there are consistent characteristics pertinent to strategies for hematophagy including: locating and accessing a suitable blood-containing host, penetrating the integument and blood circulatory system, counteracting the hemostatic system of the host, taking a blood-meal and escaping detection so as to survive to digest the meal.\nTo obtain blood, hematophagous ectoparasites cut or pierce the epithelium and disrupt the integrity of the circulatory system of the host. During this process of gaining access to blood and feeding on blood, hematophages may inject substances into the host such substances being brought to the symbiotic relationship by the parasite. These injected substances may be either produced by the parasite and assist in blood-feeding or present in the hematophage as a residue of previous feeding.\nBlood-feeding behavior results in numerous new infections in humans per year and many more in animals. Among the most widely known of these infections are cases of malaria resulting from mosquito bites. While Yellow fever, Chagas disease, Dengue fever, Ebola, and other diseases are transmitted by tropical insects, human arthropod-borne infections also occur in non-tropical regions. Infectious diseases such as Borrelia, Babesia, Rickettsia, West Nile virus, and various types of encephalitis infections are found in cold weather climates in the Northern Hemisphere. Similarly, many non-human species are infected by hematophagous parasites and may serve as alternate hosts or reservoirs for human pathogens.\nAmong blood-feeding organisms, arthropods are possibly the best known. Included in this group are arachnids (ticks and mites) and insects, flies, bugs, and lice. Some hematophagous parasites require blood-meals at every stage of life, such as is the case with ticks, lice, and bedbugs, and others only as adults and in some cases only females are parasitic. Some hematophages spend their entire life on the hosts such as is the case with head lice. Other species, such as bedbugs, get onto the source of blood only to feed. The wide biological diversity of hematophagous organisms presents a challenge to detecting and controlling these pests.\nSources of blood for these parasites are vertebrates including, most importantly, terrestrial reptiles, birds, and mammals because of their association with humans. Some species of hematophages are opportunistic and will take meals from a variety of host species whereas others may preferentially or obligately feed on one host species. In many cases humans are more exposed to hematophagous organisms upon leaving shelter. Well-known blood-sucking insects, such as mosquitoes and biting flies, and ticks, live primarily outdoors. While these parasites may enter houses, primarily to feed, they usually do not infest human shelters. Typically, arthropod life cycles are more compatible with the environment outside of human dwellings and humans are opportunistically used as sources of blood.\nThere are, however, animals that opportunistically share houses with humans even in developed countries. These organisms are commonly known as pests. Many arthropod pests enter a dwelling with other pests; for example, fleas may come in on rodents. It is also possible that pests enter dwellings on their own, with the humans themselves, or companion animals such as dogs or cats carry pests into the homes. Infestations by hematophagous ectoparasites can cause annoyance, discomfort, health problems, and economic losses. For example, pests in structures other than those for human habitation, including but not limited to hen houses for avian egg production, barns housing domestic animals, mammalian breeding facilities, or other structures for continued occupation by birds or mammals, may cause reduced productivity. Once inside, some arthropods are able to establish breeding colonies in the home. Often the homeowner is not aware of such infestations. For example, house dust mite colonies exist in a high percentage of homes in the United States without the human co-inhabitants being aware of their presence. These organisms feed primarily on human skin flakes that are sloughed constantly and are a component of house dust. House dust can also serve as the main nutrient source for immature fleas, obligate blood-feeders as egg producing adults.\nPesticide use to control many hematophagous species has undergone significant changes with the virtual elimination of “residual insecticides” (Spielman, Pollack et al. 2001). Subsequently, current pest management practices have not eliminated hematophagous ectoparasites from shelters: school children still are infected by head lice (Pediculus capitas); bedbugs (Cimex lectularius) infest henhouses and are being found with increasing frequency in homes and lodging facilities in the United States; fleas (Ctenocephalides sp.) and several genera of ticks enter structures inhabited by humans, mice, dogs, cats, and other mammals alone or in combination; and companion animals of humans are frequently hosts for are variety arachnid and insect ectoparasites. Because modern integrated pest management practices tend toward more judicious use of pesticides, determining that an infestation exists is a necessary first step in controlling pests and is useful in evaluating the effectiveness of control measures taken. Currently visual inspection either directly or using various collection and magnification aids is the primary means to determine the presence of pests. Frequent publication of articles and patents describing new and better ways to detect the presence of pests attests to the fact that current methods need improvement.\nAs mentioned above, blood-sucking parasites may inject substances into the host during feeding. Compounds of parasite origin have been found to have various physiological activities including anticoagulation, platelet aggregation inhibition, pain inhibition, and anti-inflammation. Such findings have lead to published studies describing the structure and function of these compounds of possible use as pharmaceuticals. Hirudin, produced by the leech Hirudo medicinalis, has been cloned for mass production by recombinant DNA technology and is available for treating humans as “Refludan” (BERLEX Laboratories, Wayne, N.J.). This and other therapeutic applications have lead to an expanding literature base in pharmacognosy focused on proteins, especially of salivary origin, elaborated by hematophagous invertebrates to assist in blood-feeding (Valenzuela, Walker et al. 1995; Valenzuela, Charlab et al. 1998; Valenzuela, Belkaid et al. 2001). These studies have provided DNA sequences for many salivary proteins which served as useful prior art for the present invention. Salivary biomolecules are injected into the host circulatory system where the physiological effects occur. The proteins of hematophage origin are also carried by the blood into the parasite digestive system where they become part of the blood-meal and are subject to degradation by the digestive hydrolytic processes. As a result, the structure and molecular characteristics of salivary proteins may be altered during passage through the digestive system of parasites and be different from those of the salivary proteins per se or produced from recombinant genes.\nVertebrate blood is typically about 80% water. Many hematophagous invertebrates, especially ectoparasitic species, utilize diuresis mechanisms to eliminate much of the water from a blood-meal very quickly usually as feces (Quinlan, Tublitz et al. 1997). This is necessary for many reasons including improved mobility of the pest to avoid being detected and preyed upon or swatted and more efficient digestion. The very liquid feces contain compounds that have undergone little digestion during the short time between feeding and excretion. These feces will contain biomolecular markers of blood feeding activity derived from both the host blood and those associated with the hematophagous organism involved.\nVarious types of tests have proven to be reliable and cost-effective for detecting biomolecules. Those skilled in the art of detecting molecules of biological origin have developed various methods to offer rapid, easy-to-use, and specific tests with sufficient sensitivity to be useful for detecting markers for a wide variety of organisms including but not limited to bacteria, viruses, parasites, and arthropods in the environment. Many of these assays take advantage of the chemical composition and of specific binding properties of biochemicals including but not limited to effectors, inhibitors, modulators, and hormones binding with receptors; lipids, metals, proteins, hormones being specifically bound to carrier proteins; antibodies binding to cell surface binding domains; and the well-know interaction between antigens and antibodies. Thus, while detection assays are know in the art, such assays have not been successfully adapted to the detection of hematophagous ectoparasites that live outside of the body of the host by the methods and compositions of this invention."}
{"text": "An electronic apparatus, such as a portable computer or a cellular phone, comprises a connecter through which peripheral equipment is connected. The connector is housed inside the electronic apparatus and is exposed through an opening in the periphery of the housing.\nThe apparatus also comprises a printed circuit board which has a connector mounting region in which a plurality of pads are arranged. The connector comprises a plurality of connection terminals corresponding to the pads, the connection terminals being soldered to the pads.\nThere is an increasing need to make the housing of a portable electronic apparatus, such as a portable computer, thinner. As a way of doing this, Jpn. Pat. Appln. KOKAI Publication No. 2005-268521, for example, discloses an electronic apparatus in which a conductive adhesive is printed on the inside of the housing to form a circuit pattern integrated with the housing, and connection terminals are bonded to the lands of the pattern.\nWith this arrangement, the circuit pattern and a connector are electrically connected by the conductive adhesive forming the circuit pattern, and the connector is fixed at a predetermined position on the housing. This eliminates the need for a printed circuit board, saves space and makes the housing thinner.\nJpn. Pat. Appln. KOKAI Publication No. 2005-268521 also discloses that the circuit pattern formed by the conductive adhesive takes the form of thin lines. This structure makes it difficult to ensure a sufficient connection area between the circuit pattern and the connection terminals of a connector.\nThus, sufficient adhesive strength between the connector and housing may not be achieved, impairing the soundness of the electrical and mechanical joining of the connector."}
{"text": "Various gas sterilization methods have been investigated in the past. Methods using ethylene oxide and other disinfecting gases are widely used for sterilizing a wide range of medical products from pharmaceutical preparations to surgical instruments. Irradiation alone or together with disinfecting gases has also been investigated, as summarized by Russell, A., The Destruction of Bacterial Spores, New York: Academic Press (1982).\nA sterilizing method must effectively render all microbial organisms non-viable without damage to the article or goods being sterilized and its packaging. However, many disinfecting gases which meet this criteria, such as ethylene oxide and irradiation methods, have been recognized to expose workers and the environment to safety hazards. Recent legislation has been severely restricting the amount of hazardous gases such as ethylene oxide (a suspected carcinogen) in the working environment, or the use of any system or method which produces toxic residues or exhaust products. This has been presenting a major crisis in hospitals and other areas of the health industry."}
{"text": "Modern network units, by which one principally means switches and routers, have reached a substantial degree of sophistication and can be realized in a considerable variety of architectures. It is not intended to limit the invention to any particular architecture. Broadly, units of this nature can switch packets in accordance with either media access control address data (layer 2 in the OSI Model) or network address data (layer 3 in the OSI Model) and/or possibly information in higher layers. For this purpose a received packet may be temporarily stored, either in temporary storage associated with a respective ingress port, or in central memory which may have memory space or buffers permanently or temporarily associated with respective ports, in queues which may be constituted by the packets themselves along with status data or which may be constituted by pointers to locations in memory. Header information from the packets is subject to a look-up in order to obtain, with recourse to a database, forwarding information for the packet. This forwarding information may be subject to post-processing in accordance with a variety of rules, such as VLAN matching, quality of service, spanning tree rules and others which are important but have no direct bearing on the present invention. The intended result is always to determine, on the assumption that the packet is to be forwarded from the unit, which port or (in the case of a multicast packet) ports are to be selected for the forwarding of the packet from that unit. A switching engine typically is controlled by a port mask, obtained as a result of the look-up. The port mask identifies, directly or indirectly, the ports of the unit by number and (for a unicast packet) the port from which the packet should be forwarded. Generally speaking, ‘switches’ or ‘bridges’ are terms applied to units that switch in layer 2 and ‘routers’ is a term generally employed for switching in layer 3, since media access control data is employed to determine a device to which the packets should be forwarded whereas network address data relates to a network to which a packet should be forwarded. However, usage of these terms is not exclusive, because units capable of both bridging and routing are known. Accordingly, in the present application ‘network unit’ is employed to refer to a device which performs the forwarding of data to a selected port or ports having regard to address data in the packet and, optionally, other data which may relate to the type of packet and/or other data which may affect the forwarding decision.\nCascade Systems\nNetwork units are made with a fixed number of ports which are employed for the reception of packets from and the forwarding of packets to the external network. It is generally now considered convenient to manufacture switches in a comparatively small number of individual ‘sizes’, in terms of the number of ports, and to provide in effect a network unit with a much larger number of ports by means of a ‘cascade’ system, wherein a multiplicity of multiport units are connected and organized so that they appear to the external network as a single switching or routing entity.\nVarious desirable features in modern network practice, the existence of a variety of forwarding rules, and the practical importance of employing ‘resilience’, has led to the development of sophisticated cascade systems.\n‘Resilience’, as previously mentioned, is the feature characteristic of a system which can tolerate the powering-down or failure of a unit or a connecting link and which will maintain the ‘cascade’ or ‘stack’ in operation as far as the other units in the cascade are concerned. A good example is the system described in GB-A-2365718. That document describes a cascade system in which the network units in a stack are coupled by a special connector, known as a T-piece, which can provide continuity of a data path for packets in a ring connecting the units notwithstanding the powering-down or failure of a unit to which it is coupled. For that purpose the connectors contain a system of multiplexers which are under the control of control logic which determine the status of the various network units by means of the exchange of control frames with the network units.\nAlthough a system as described in that document is satisfactory in operation, the reliance on special connectors imposes a hardware overhead and, moreover, constrains the cascade system to have the configuration of a ring. In practice also, the network units have to be in reasonably close proximity. Furthermore, the traffic volume for packets between the units is limited by the capacity of the data path through the connectors.\nIn a network unit which is in a cascade connection, a look-up has in general three possible results. The address may be ‘unknown’ in that there is no entry in the database which can determine the port or ports from which the packets should be forwarded. Such a result may have the consequence of ‘flooding’: the packet is transmitted out of all the ports of the switch, or possibly out of a group associated with a particular virtual segment of the network, if partitioning in the form of virtual LANs is used. A response from a destination conforming to the destination address may be established by means of an address resolution protocol (ARP). Second, the destination may be by way of a port on the respective unit, that is to say by way of a ‘local’ port. In such circumstances the packet can be forwarded out of the relevant port without recourse to the cascade connection. Thirdly, the destination may be by way of a port on another switch in the cascade, in which case the port number in the forwarding data would usually be a cascade port or possibly a ‘logical’ port connected to a group of cascade ports.\nIt is customary for look-up databases to include an address learning facility. Typically the source address data of an incoming packet is checked against entries in the database and if the source address is unknown it can be ‘learnt’ against the port by which it is received. It is also possible to insert entries in such databases.\nOne object of the present invention is to permit network units to be in a cascade system which may have a wide variety of configurations, and in particular a general mesh which provides a multiplicity of potential paths between units.\nAdditionally or alternatively, a further object of the invention is to permit the network units which constitute the cascaded system to be substantially widely spaced, for example in different buildings. Such a facility would substantially extend the versatility of cascade systems.\nA further object of the invention is to control the forwarding of packets within the cascade in a manner which tolerates and indeed facilitates a variety of configurations including meshes.\nA yet further object of the invention is to facilitate trunking between different units in the cascade connection; i.e. the provision of a multiplicity of parallel connections between units to increase the traffic capacity between them."}
{"text": "When a user wants to use an electronic product, such as a handset having an Internet-access function, to gain wireless access to the Internet, a long string of codes must be manually inputted, such as Wired Equivalent Privacy keys (WEP Keys), to register with an access point so as to be able to use the access point to activate the Internet-access function in an encrypted and protected environment. The purpose of inputting the codes is to activate the encryption mechanism so that hacker attacks or data interception can be avoided when the user accesses the Internet wirelessly.\nWhen the user holding the electronic product moves to an area covered by another access point, the user must register with the new access point. That is, another long string of codes must be inputted in order to activate the Internet-access function through the access point. Therefore, the conventional security authentication mechanism not only requires the user to be familiar with the codes of every access point, the user also needs to input the codes manually, thereby resulting in great inconvenience in use.\nFor a moving user, he/she may cross areas covered by several access points during the moving process, and may need to input a long string of specific codes upon switching to a new access point. Therefore, some access point administrators set the codes of the access points administered thereby to be the same or even disable the code settings of such access points so that users no longer need to input specific codes in order to minimize user inconvenience. However, regardless of whether the encryption mechanism is disabled or the access points are set to have the same codes, these are passive methods of handling the above problem, and may result in hacker attacks which will expose users to danger when accessing the Internet wirelessly."}
{"text": "1. Field\nThe present invention generally relates to controlling connection among a plurality of transmission terminals.\n2. Background\nWith the need for reducing costs or times associated with business trips, more companies are moving towards videoconference systems to have teleconference or videoconference among remotely located offices via a communication network such as the Internet. The videoconference systems allow transmission of contents data such as image data and/or sound data among a plurality of transmission terminals that are remotely located from one another through the communication network to facilitate communication among the plurality of transmission terminals.\nTo have a videoconference among the plurality of transmission terminals, a user is required to previously register information regarding participants, who will be participating in the videoconference before the videoconference starts. In case a user who is not previously registered decides to participate in that videoconference after the videoconference starts, the unregistered user is not allowed to participate unless the unregistered user is registered through the registration process. The registration process usually requires additional time such that it has been inconvenient for the unregistered user to participate in the videoconference after the videoconference starts.\nJapanese Patent Application Publication No. H08-256145 describes a conference manager, which manages a permitted user list that registers identification of a user who is allowed to participate in a conference without the registration process, for each of conferences. When a request for participating in a specific conference is received from a user, the conference manager determines whether identification of the user is registered in the permitted user list prepared for the specific conference, and allows the user to participate in the specific conference if the user is registered in the permitted user list."}
{"text": "This invention relates to a shifting apparatus, and more particularly to a paper shifting apparatus designed to automatically replace a continuous toilet-seat covering paper on toilet seats.\nThe development of modern toilet fixtures constitutes one of the most remarkable improvements in the living conditions of human beings. Unfortunately, so far as the inventor knows, no steps have been taken to improve the sanitary conditions of toilet seats. Today the only sanitation precautions available are the cleaning operations performed by the users with toilet paper and manually applied cleansers. In order to solve this sanitary problem, a paper toilet-seat covering is developed for personal use. However, as the known paper toilet-seat cover is normally prepared in a roll and placed at a location close to a toilet, to use it, a user has to draw a piece of the paper therefrom and spread it on the toilet seat. This known paper toilet-seat cover is rather inconvenient and often causes stoppage of the toilet bowl after the used paper toilet-seat cover has been dropped into the toilet bowl."}
{"text": "Surfacing machines, such as wood planers and jointers, etc., often make use of at least one powered, rotatable cutter head. The cutter head is typically supplied with a number of axial knives presenting sharpened edges along the outward periphery of the cutter head for removing material passed over the rotating head on a workpiece support surface. In wood planers, often two or more such cutter heads are used to simultaneously perform surfacing operations on opposing surfaces of the wooden workpiece.\nThe process involved in mounting the cutter knives to rotatable cutter heads is rather complicated and time consuming. It is therefore desirable to periodically recondition or \"joint\" the knives on the cutter head rather than removing the knives and sharpening each separately.\nProper equipment used in the jointing operation and correct selection of jointing stones are major assets in accomplishing a proper jointing operation.\nA skilled worker may have adequate training and expertise to accomplish a quality job, but may well be frustrated in doing so because of equipment shortcomings. For example, a soft jointing stone is almost always preferred over a hard, dense stone. But the soft stones cannot be used due to the pounding action of the rotating knives against the stone if the stone support \"flutters.\" This is a typical situation with most, if not all, present \"dovetail\" jointer stone support assemblies. A worn or loose fit in the stone support allows the stone support to jump or \"flutter\" as the knives rotate into engagement with the stone. Soft stones are easily broken in this situation so the less desirable hard stones must be selected.\nAnother result of the stone support flutter is that an undesirable high or outwardly projecting \"heel\" is created behind the cutting edge of the knives. This excessive \"heel\" does not assist the knives in the cutting or scraping operation, but instead rubs over the surface of the workpiece, often causing burns or, at best, a damaged textured surface. Nonetheless, the worker must select a hard stone for the jointing process since a soft stone would be surely broken by the fluttering action of the support and contact with the surfacer knives.\nA tight or securely mounted stone holder will facilitate the use of softer stones which remove many times more material in a shorter time. However, the presently available jointer devices cannot be effectively \"tightened\" to provide firm support for the stone without sacrificing freedom of movement for the stone across the cutter head during the jointing operation. Dovetail stone mounts will not slide freely when secured in a manner sufficient to reduce or eliminate stone flutter. Older, worn jointing devices simply cannot be adjusted tight enough to accomplish the task, with the inherent sacrifice of mobility of the stone across the cutter head.\nThe above problems have been recognized to limited degree by the below-referenced United States patents. However, it is not believed that any of the known references show or suggest the solution embodied in the present invention as set forth below.\nU.S. Pat. No. 4,495,734 granted to Rauch in 1985 discloses a grinding arrangement for chopping cutters. The grinding arrangement makes use of a tool mounting bar and guide slidably engaged thereon for supporting the grinding stone. A wheel actuator arrangement is provided for driving and guiding the stone to and fro along the axial length of the cutter head. This arrangement, while making use of wheels in the drive arrangement, still makes use of a singular bar and guide way arrangement for mounting the grinder.\nU.S. Pat. No. 3,374,699 to Schmermund discloses a sharpening arrangement for mounting cutter heads and for sharpening the cutter heads. The apparatus includes a support frame for receiving and mounting the cutter head for rotation on an axis that is skewed in relation to a grinding tool support. The grinding tool is movably supported on a pair of guide rods for motion along the cutter head. A stone support is mounted on the rod-mounted carriage by a dovetail arrangement to facilitate radial adjustment of the stone against the cutter head to be sharpened.\nThe rods and bearings mounting the carriage and stone in the Schmermund device may function to eliminate some of the problems present in a standard dovetail guide arranged but many of these advantages may be offset by use of the dovetail stone support arrangement. It may therefore be concluded that Schmermund, while providing adequate longitudinal support for the jointing stone, teaches away from the solution presently disclosed in this application by using a dovetail stone mounting arrangement for radial adjustment relative to the blade or cutter head to be sharpened.\nU.S. Pat. No. 2,476,177 to Bloom et al. issued Jul. 12, 1949, discloses a forage harvester knife sharpener. This device again makes use of an elongated bar mounting a stone holder for adjustment toward or away from the cutters on a drum. The entire mechanism is pivotal about a shaft that extends across the width of the harvester cutting head.\nU.S. Pat. No. 993,398 to Osborne granted May 30, 1911, discloses a grinder for planer knives also making use of a bar slidably mounting a jointing stone support.\nU.S. Pat. No. 267,579 to M. W. Palmer granted Nov. 14, 1882, discloses a planer knife grinder in which a grinding disc is mounted for rotation with an end-supported shaft adjacent a planer knife head. The powered grinding disc can be moved along the length of the rotatable single shaft by an elongated actuator extending to one side of the planer. The Palmer device is intended to sharpen each knife separately and is not a jointer cannot be used effectively with the planer in operation. A somewhat similar apparatus is disclosed in U.S. Pat. No. 2,620,606 to Dvorak, issued Dec. 9, 1952. Dvorak discloses a sharpening attachment for surfaces in which a grinding wheel is mounted by a pivoted frame to a shaft for sliding movement along the cutter head. The cutter head cannot be operated during the sharpening operation.\nOf the above references, none are believed to show or suggest a solution to the problem of securing a jointing stone while enabling relatively free precision motion of the stone in relation to a rotating cutter head as described and claimed below."}
{"text": "1. Field of the Invention\nThe present disclosure relates to a secondary battery and a battery pack, and particularly relates to a secondary battery including a current cut-off device, and a battery pack configured by connecting a plurality of secondary batteries of this type.\n2. Description of Related Art\nJapanese Patent Application Publication No. 2011-258550 discloses a secondary battery having a configuration that deforms a short-circuiting member by increase in inner pressure of a case so as to short-circuit a positive electrode terminal and a negative electrode terminal to fuse a fuse portion by an overcurrent, thereby cutting off an electric connection between the positive electrode terminal and an electrode body.\nIn the configuration as described in JP 2011-258550 A, the short-circuiting member is formed to be thin enough to be easily inverted by the pressure. Hence, the short-circuiting member might be fused before the fuse portion is fused due to a large current flowing in a moment at the time of the short circuit.\nIf the short-circuiting member is fused but the fuse portion does not work, connection between the terminal and the electrode body is maintained, so that increase in inner pressure of the case is continued. Consequently, the secondary battery might be damaged.\nOn the other hand, if the short-circuiting member is fused and the fuse portion works, the connection between the terminal and the electrode body is cut off, and no short circuit between the positive electrode terminal and the negative electrode terminal is formed, and thus a cut-off of conduction among batteries is caused by an abnormal secondary battery. Consequently, it becomes impossible to pick up electric power from normal batteries, which disables charging and discharging in the whole battery pack."}
{"text": "1. Field of the Invention.\nThe present invention relates generally to spacing devices, and more particularly, but not by way of limitation, to a spacing member for connecting a voltage regulator in a spaced relationship to a housing of an alternator so as to enhance heat transfer away from the voltage regulator by creating an effective air flow path.\n2. Brief Description of the Related Art.\nAlternators for automobiles are provided with a housing in which a rotor coil is supported. The engine of the automobile is utilized to rotate the rotor coil to generate a rotating magnetic field that induces a current in stator windings disposed near the rotor coil. The speed at which the rotor coil rotates and the intensity of the rotor coil's magnetic field determine the voltage and current of the electricity generated by the alternator. To control the intensity of the magnetic field generated by the rotor coil so that the output voltage of the alternator remains within predefined limits, a voltage regulator is generally employed with an alternator.\nThe positioning of voltage regulators relative to the alternator housing is varied with different makes and models of alternators. For example, with some alternators the voltage regulator is mounted a distance from the alternator. In other alternators the voltage regulator is mounted within the alternator housing. Finally, in some makes and models, in particular Ford IAR 3G alternators, the alternator housing is provided with an external mounting surface to which the voltage regulator is mounted.\nOverheating of the voltage regulator is a primary cause of alternator failure. The voltage regulator is subjected to both external heat produced from the automobile engine and internal heat produced by the flow of electric current through the voltage regulator itself. And while many voltage regulators are provided with a heat sink to transfer the heat away from the voltage regulator, the manner in which the voltage regulators are designed to be mounted to the alternator housing does not permit adequate air circulation across the heat sink and thereby prevents heat from being effectively transferred away from the voltage regulator.\nTo this end, a need exists for a voltage regulator spacing member which permits the voltage regulator in certain models of alternators to be spaced from the alternator housing so as to increase air flow through the voltage regulator heat sink and thus increase heat transfer away from the voltage regulator. It is to such an apparatus that the present invention is directed."}
{"text": "The present invention relates generally to a system and method for motor control, and more particularly, to a subunit for a motor control center which connects motor control components to supply power, and a method of manufacture thereof. The system and method described herein provide for connection of the supply power to the motor control components after full installation of the motor control center subunit into the motor control center.\nA motor control center is a multi-compartment steel enclosure with a bus system to distribute electrical power, on a common bus system, to a plurality of individual motor control units mountable within the compartments. The individual motor control units are commonly referred to as “buckets” and are typically constructed to be removable, pull-out units that have or are installed behind individual sealed doors on the motor control center enclosure. These buckets may contain various motor control and motor protection components such as motor controllers, starters, contactor assemblies, overload relays, circuit breakers, motor circuit protectors, various disconnects, and similar devices for operating high voltage motors. The buckets connect to the supply power lines of the motor control center and conduct supply power to the line side of the motor control devices, for operation of motors. Motor control centers are most often used in factories and industrial facilities which utilize high power electrical motors, pumps, and other loads.\nTypically, when installing or removing motor control center buckets, the power supply lines are connected or interrupted. To remove such a bucket, a deadfront door of the bucket or of the motor control center is opened and an operator manually pulls on the bucket to separate the primary disconnects, or “stabs,” from the bus system, thereby disconnecting power supply. Installation of a bucket is accomplished in a similar manner, wherein the operator manually pushes the bucket into a compartment of the motor control center to engage the bucket stabs with the bus system, and thus connect the system to supply power. The line connections or stabs may be difficult to maneuver manually when an operator is supporting the entire bucket or when the stabs are not visible.\nAttempts have been made to improve upon the manual installation and disconnection of motor control center buckets and supply power connections from live supply power lines, risers, and/or a vertical bus of a motor control center. Other systems have employed pivotable handles inside the buckets to pivot line connectors to and from supply lines. However, many of these systems require that the bucket or compartment door be open to manipulate the handles and line stabs.\nIt would therefore be desirable to design a motor control center bucket assembly that overcomes the aforementioned drawbacks. Thus, it would be desirable to provide for remote connection or disconnection of the line stabs of a bucket to the power supply lines or bus of a motor control center from a distance. In the event of an arc or arc flash, any heated gas, flame, and/or the arc itself should preferably be contained behind the bucket compartment door or “deadfront.”"}
{"text": "The present invention relates generally to molded articles which may be blow molded such as a blow-molded ladder or a blow-molded staircase for a swimming pool.\nSwimming pool ladders or staircases and various other articles have been made from a plastic material by a process known as “blow-molding,” as exemplified in commonly assigned U.S. Pat. No. 4,023,647 to Confer, commonly assigned U.S. Pat. No. 4,067,614 to Confer et al, and U.S. Pat. No. 4,166,833 to Schurman, all of which patents are hereby incorporated herein by reference. In the process of blow molding, a hollow thin-walled structure is formed having an exterior surface separated from an interior surface by the thin wall, and the hollow structure contains a fluid, in most cases air and/or water.\nPool Ladders with Vandal Guards\nBy way of background, swimming pool ladders of an A-frame type are positioned straddling a pool wall with the pool side having a ladder for entering and leaving the water (a.k.a., in-pool ladder) and an outside ladder for climbing over the outside of the pool wall (a.k.a., staircase ladder). In the past certain ladders of the foregoing type had relatively unwieldy structures for placing the outside ladder in an inaccessible position so as to prevent unauthorized entry into the pool. Also, insofar as known, ladders of the foregoing type, when disassembled, were not of a size which would fit into a box which was within the dimensions acceptable to commercial shippers. Other types of pool ladders were relatively complicated and difficult to assemble.\nIn U.S. Pat. No. 6,880,674; St-Hilaire wrote, “A door ladder assembly for use with the [sic] an above ground swimming pool, the assembly comprising a ladder having a plurality of steps, a door with a first side hingedly connected to one side of the ladder and arranged to hingedly move between open and closed positions, a handle being located at the top portion of the door, and a door lock to lock the door in a closed position, the door being provided with a lock operator to unlock the lock, the lock operator being located proximate the handle. The arrangement provides for easy access for an adult, while preventing a child from having access to the unlocking mechanism.” A prior version of a door gate for a pool ladder was disclosed in U.S. Pat. No. 3,225,863 to Ludlow and U.S. Pat. No. 3,968,857 to Bryan. Each of these prior pool ladder barriers are doors that swing on a hinge and are positioned over the A-frame ladder's staircase ladder. These barriers are variations of conventional vandal guards used in other industries.\nAn alternative vandal guard is disclosed in U.S. Pat. No. 4,579,197 to Spurling. Spurling wrote, “A ladder shield for use in preventing toddlers and the like from climbing the steps of a ladder, such as an above-ground swimming pool ladder, when the ladder or the device to which the ladder is coupled is unattended is disclosed. The ladder shield comprises a rigid yet slightly flexible shell of plastic which is sized and shaped to fit around the front and sides of the ladder so as to block access to the steps on the ladder. The ladder shield is removably secured to the ladder by a removable locking bar which when inserted extends through a pair of holes in the side walls of the shield. The locking bar is secured in place by a lock. When not being used to prevent access to the steps of the ladder, the ladder shield may be used as a mini-foot wash.”\nLadders with a Tambour\nA tambour device is defined at www.dictionary.com as “a flexible shutter used . . . in place of a door, composed of a number of closely set wood strips attached to a piece of cloth, the whole sliding in grooves along the sides or at the top and bottom.” Applicant conducted a search to determine if any ladder was associated with a tambour device. The closest reference, not a relevant reference, was U.S. Pat. No. 5,046,582 to Albrecht for a foldable ladder combination with truck cargo carrier. Albrecht wrote, “Commercial truckers frequently have situations arise where they require personal access to the elevated bed of the truck's cargo carrier. The cargo carrier may be mounted on the truck frame itself as in the case of a so-called bob-tail truck; or, the cargo carrier may be a trailer pulled by a truck tractor. In either case, the conventional truck cargo carrier includes an elongated approximately horizontal cargo bed elevated about four to five feet above ground level, and an elongated structural cargo enclosure extending upwardly from, and substantially covering the cargo bed. The cargo bed has a rear end from which cargo is loaded into and unloaded from the cargo carrier, and the structural cargo enclosure has at its rear end right and left side edges in the form of vertical posts extending upward from adjacent the rear end of the cargo bed, and typically has either a single “roll-up” tambour rear door which is engaged in slots disposed at the inboard lateral faces of the posts, or has a pair of swinging doors hinged to such vertical posts, for opening and closing the rear end of the structural cargo enclosure to provide loading and unloading rear access to the structural cargo enclosure and cargo bed.”\nLocking Mechanisms for Blow Molded Materials\nWalter, in U.S. Pat. No. 1,911,539, Doernemann, in U.S. Pat. No. 5,536,111, Troester, in U.S. Pat. No. 5,013,508, and Tokunaga, in U.S. Pat. No. 3,759,043, all disclose a three-dimensional polymeric structure having a male and a female interlocking component. Not one of these disclosures has both interlocking components on a single side of the structure. Moreover, these references disclose both interlocking components contacting the edges of the sides, which allows easier forces to separate the components. As such, none of these references discloses an apparatus that forms at least a male and a female interlocking component on a single side, and preferably not contacting an edge of the side, of the structure.\nIn commonly assigned U.S. Pat. No. 7,234,780; Lipniarski wrote, “A device and method for forming an indented female connector where the dimension of the indented female connector provides for removing a male molding component without damaging a set of internal ribs formed inside the female connector. A structure comprising an indented female blow-molded connector having a first wall having a top surface and a thickness, with the first wall surrounding an opening. At least one inner wall extends substantially perpendicular to the first wall and borders the opening. At least one rib is formed in the inner wall. A bottom wall is joined to the inner wall and the distance from the rib to the top surface is greater than the thickness of the first wall. The structure further comprising a male connector having a contact edge that, when positioned inside the indented female blow-molded connector, engage the at least one rib, such that the male connector and indented female blow-molded connector are releaseably joinable with one another. Structures can be made that use the indented female blow molded connectors and male connectors including braced stools having three or four legs.” This connection system requires the male component be inserted immediately into a female locking area (a.k.a., an immediate locking system), which is difficult for easy and secure assembly.\nSince most blow molded polymeric materials have an immediate locking system, most blow molded polymeric materials having an immediate locking system also require additional interconnection devices be used to secure the assembly. Those additional interconnection devices include and are not limited to rods, screws and other conventional interconnection devices. See U.S. Pat. No. 6,769,513.\nThe present invention is directed to a blow molded plastic object having a female interconnection system comprising a receiving area and a locking area. We are aware of one patent that discloses a similar a female interconnection system. That patent is commonly assigned U.S. Pat. No. 6,190,600; which is hereby incorporated by reference. In that patent, Lipniarski wrote, “The present invention relates to an apparatus that forms a polymeric structure. The apparatus has a mold plate, a circular trench, a male and female interconnect, a vertical gap filler, and a cylinder. The mold plate has a predetermined geometric design. The circular trench is in the mold plate and the difference between the outer diameter and the inner diameter of the circular trench is a distance D. The male interconnect and the female interconnect are interspaced between each other. Each interconnect rotates within the circular trench and has a base and an interconnection portion. Each base has a width D and slidably mates to the mold plate. Each interconnect portion has a minimum width W which is less than D. The male interconnect portion protrudes downwards relatively from its base to a distance P, likewise the female interconnect portion protrudes upwards relatively from its base to a distance P. The vertical gap filler aligns with the circular trench and has a foundation and an extension. The extension has a width greater than W, a height greater than 2P, and protrudes upwards and downwards from its foundation. The cylinder rotates each interconnect to a predetermined position and raises the vertical gap filler so the foundation connects to the mold plate before any polymeric material is applied to the apparatus to form the polymeric structure. Once the polymeric structure is formed, the cylinder lowers the vertical gap filler until the extension is below the mold plate a distance greater than P and rotates each interconnect so the polymeric structure can be removed from the apparatus.”\nThe claimed device in U.S. Pat. No. 6,190,600 is a mold for creating a first blow molded object that had a first male structure having a first extension and a first block and a first female structure having first receiving area and a first inhibiting area; and a second blow molded object that had a second male structure having a second extension and a second block and a second female structure having second receiving area and a second inhibiting area. Thereby the first blow molded object and the second blow molded object could have the first male structure inserted into the second receiving area and the second male structure inserted into the first receiving area. Rotating the first blow molded object in relation to the second blow molded object so the first male structure slides into the second inhibiting area and the second male structure slides into the first inhibiting area. The first and second inhibiting areas have no blow molded locking mechanism to secure the respective male structures in the inhibiting areas. Instead the inhibiting areas rely on the walls (which surround [that means a back wall] the male structure except where the male interconnect slides from the receiving area to the inhibiting area and the male extension that protrudes through an opening in the female inhibiting area) to provide a friction fit to inhibit the male structure from releasing itself from the female structure. The lack of a locking mechanism is acceptable for wave dispersion systems but not for ladders which require a more secure system. The present invention solves this problem for blow molded structures that require a secure locking mechanism."}
{"text": "α-Synuclein is a neuronal protein which originally has been associated with neuronal plasticity during Zebra finch song learning. Although its role at the molecular level is at present largely elusive it appears to have lipid bi-layer (or membrane) with binding properties important for preserving proper transport of neurotransmitter vesicles to the axonal ends of neurons presumably to ensure proper signalling at the synapse. Apart from its physiological role in brain cells, human α-synuclein also possesses pathological features that underlies a plethora of neurodegenerative diseases including Parkinson's disease, diffuse Lewy body disease, traumatic brain injury, amyotrophic lateral sclerosis, Niemann-Pick disease, Hallervorden-Spatz syndrome, Down syndrome, neuroaxonal dystrophy, multiple system atrophy and Alzheimer's disease. These neurological disorders are characterised by the presence of insoluble α-synuclein polymers or aggregates usually residing within neuronal cells, although in the case of Alzheimer's disease α-synuclein (or proteolytic fragments thereof) constitutes the non-amyloid component of extracellular “amyloid-β plaques”. It is widely believed that the amyloidogenic properties α-synuclein disrupt cellular integrity leading to dysfunctioning or death of affected neurons resulting in cognitive and/or motoric decline as it is found in patients suffering from such diseases. The aggregation of α-synuclein is at present very poorly defined, but constitutes most likely a multi-step process wherein self-polymerization of α-synuclein into insoluble aggregates is preceded by the formation of soluble protofibrils of α-synuclein monomers. Self-association may be triggered by the formation of alternative conformations of α-synuclein monomers with high propensity to polymerize. Several studies using neuronal cell lines or whole animals have shown that formation of reactive oxygen species (hereinafter abbreviated as ROS) appear to stimulate noxious α-synuclein amyloidogenesis. For instance paraquat (an agent stimulating ROS formation within the cell) has been recognized as a stimulator of α-synuclein aggregation. Like in animals, exposure to paraquat is believed to induce the formation of synuclein inclusions, and consequently neurodegeneration, especially of dopaminergic neurons in humans. Dopaminergic neurons appear to be particularly sensitive because the concurrent dopamine metabolism may on the one hand contribute significantly to the oxidative stress load but may on the other hand result in kinetic stabilisation of highly toxic protofibrillar α-synuclein species by dopamine (or its metabolic derivatives). Parkinson's disease is characterised by a selective loss of dopaminergic substantia nigra cells and therefore treatment of animals (or neuronal cells) with paraquat is a common well-accepted experimental set-up for studying synucleopathies, in particular Parkinson's disease.\nApart from ROS, mutations in the coding region of the α-synuclein gene have also been identified as stimulators of self-polymerization resulting in early disease onset as it is observed in families afflicted by such mutations. Finally, increased expression of α-synuclein also promotes early disease onset as evidenced by a duplication or triplication of the α-synuclein gene in the genome of some individuals. The molecular mechanism by which α-synuclein self-association triggers cellular degeneration is at present largely unknown. Although it has been speculated that insoluble aggregates affect cellular integrity, it has recently been suggested that soluble protofibrillar intermediates of the aggregation process are particularly toxic for the cell as opposed to mature insoluble fibrils which may be inert end-products or may even serve as cytoprotective reservoirs of otherwise harmful soluble species. Therapeutic attempts to inhibit formation of insoluble aggregates may therefore be conceptually wrong, possibly even promoting disease progress.\nWhile the identification of pathological α-synuclein mutations unequivocally revealed a causative factor of a plethora of neurodegenerative disorders, treatments ensuring suppression of toxic α-synuclein amyloidogenesis are presently not available. Only symptomatic treatments of Parkinson's disease exist, which aim e.g. at increasing dopamine levels in order to replenish its lowered level due to degeneration of dopaminergic neurons, for instance by administrating L-DOPA or inhibitors of dopamine breakdown. Although such treatments suppress disease symptoms to some extent, they are only temporarily effective and certainly do not slow down ongoing neuronal degeneration.\nThus there is a need in the art for designing new drugs for therapeutic treatments that target the underlying molecular mechanism of α-synuclein related pathologies in order to reduce neuronal cell death and/or degeneration.\nWO 99/51584 discloses 5-piperazinyl-1,2,4-thiadiazoles as inhibitors of proton pump H+/K+-ATPase and therefor useful in the treatment of peptic ulcer. However these compounds are not suggested for use in the prevention or treatment of neuro-degenerative disorders."}
{"text": "A cable splice such as found in a telecommunications cable is ordinarily housed within a protective cover known as a cable splice enclosure. Such enclosures often have cylindrical covers with one or more longitudinal joints and circular end plates that surround incoming and outgoing cables. Seals between the plates and the cables, and between the plates and the covers protect the splice from contamination. Usually, two rigid metal bars, known as torque bars, extend longitudinally on diametrically opposite sides of the enclosure, between the oppositely facing end plates. The splice must be kept within the bounds of the enclosure, as defined by the two torque bars and the two end plates, to ensure the cylindrical covers can be fastened to the end plates.\nWhen a cable splice is made, individual wires in the cable are separated from bundles of wires. Typically bundled wires are densely packed within the cable and therefore occupy only a minimum amount of space. When the individual wires are separated however, a greater amount of space is required. In addition, cables are typically connected together using crimped connectors which also occupy a considerable amount of space. Thus the splice often occupies a large amount of space within the enclosure.\nThe splice can also be quite heavy, tending to hang downwards below the bounds of the enclosure. The torque bars typically extend on either side of the splice i.e. the bars are spaced generally horizontally apart, and therefore provide no direct support for the splice. It is, however, desirable to provide support for the splice to confine it within the bounds of the enclosure so that individual wires do not rub against the enclosure or hang outside the enclosure.\nUsually, after a splice is made, individual wires are tied into bundles using waxed string. This enables the individual wires to be tightly drawn together to occupy less space within the enclosure. After the bundles are tied, any of several methods is employed to support and protect the splice. One method involves the use of a non-adhesive, approximately 11/4inch wide woven cloth tape known as \"U.G. tape\". The tape is tied to a torque bar, looped under the splice and wrapped or tied around the other torque bar. This is repeated several times along the length of the splice to create a cradle formed by a series of loops extending between the two torque bards. Another method of making a cradle involves a similar tying and looping process but instead of the U.G. tape, a one-inch fibre glass tape is used. After a cradle has been created, the entire length of the splice is further wrapped helically in a fine pitch using an approximately three-inch wide non-adhesive cloth tape known as muslin tape. The splice is wrapped tightly to ensure the bundles of wires are drawn radially inwardly within the bounds of the enclosure. The enclosure covers are then installed and the repair operation is completed.\nThe use of the above methods for wrapping a cable splice presents problems to repairmen who are later required to open the splice for maintenance or other work. In order to gain access to any particular wire pair within the splice, the repairman must unwrap the muslin tape to expose the entire splice while the U.G. tape or fibre glass tape remains intact, unless access to a lower portion of the splice is necessary, in which case the cradle is cut. The repair or maintenance work can then be performed, after which the repairman is required to recreate the cradle if the previous cradle was cut and is required to rewrap the splice with muslin tape. This can be very time consuming.\nOften, cable splice enclosures are located in manholes or on above-ground cables rendering them susceptible to water contamination while the enclosure covers are removed. It will be appreciated that water can be retained within the muslin tape. Should the enclosure covers be installed while the muslin tape is wet, the tape is susceptible to rotting and the cable splices are susceptible to corrosion. Thus, should water come in contact with the muslin tape, the tape must be dried before the enclosure covers are installed."}
{"text": "Early attempts to modify or enhance the color of one's eyes utilized colored contact lenses with a simple solidly colored area that covered the iris portion of the eye. However, contact lenses with this type of opaque coloring imparted a very unnatural appearance. Other types of colored contact lenses were developed, such as Wichterle, U.S. Pat. No. 3,679,504, which discloses an opaque lens having an iris of more than a single color artistically drawn or photographically reproduced. However, such lenses did not look natural and as such never achieved commercial success. Other attempts to produce an opaque lens with a natural appearance are disclosed in. U.S. Pat. No. 3,536,386, (Spivak); U.S. Pat. No. 3,712,718 (LeGrand), U.S. Pat. No. 4,460,523 (Neefe), U.S. Pat. No. 4,719,657 (Bawa), U.S. Pat. No. 4,744,647 (Meshel et al.), U.S. Pat. No. 4,634,449 (Jenkins); European Patent Publication No. 0 309 154 (Allergan) and U.K Patent Application No. 2 202 540 A (IGEL).\nCommercial success was achieved by the colored contact lens described in Knapp (in U.S. Pat. No. 4,582,402) which discloses a contact lens having, in its preferred embodiment, colored, opaque dots. The Knapp lens provides a natural appearance with a lens that is simple and inexpensive to produce, using a simple one-color printed dot pattern. Although the intermittent pattern of dots does not fully cover the iris, it provides a sufficient density of dots that a masking effect gives the appearance of a continuous color when viewed by an ordinary observer. Knapp also discloses that the printing step may be repeated one or more-times using different patterns in different colors, since upon close examination the iris is found to contain more than one color. The printed pattern need not be absolutely uniform, allowing for enhancement of the fine structure of the iris. The one-color Knapp lenses currently achieving commercial success have their dots arranged in an irregular pattern to enhance the structure of the iris. However, neither the Knapp commercial lenses, nor the Knapp patent disclose or suggest how one would arrange a pattern of dots having more than one color to achieve a more natural appearance.\nVarious efforts have been made to improve on the Knapp lens. U.S. Pat. No. 5,414,477 to Jahnke discloses the application of the intermittent ink pattern in two or more portions of distinct shades of colorant to provide a more natural appearance.\nOther attempts to create a more natural appearing lens include U.S. Pat. No. 5,120,121 to Rawlings, which discloses a cluster of interconnecting lines radiating from the periphery of the pupil portion to the periphery of the iris portion. Further, European Patent No. 0 472 496 A2 shows a contact lens having a pattern of lines that attempts to replicate the lines found in the iris.\nDespite these efforts, the contact lens industry continues to seek a low-cost, colored lens that can enhance or modify the eye color, while providing the depth and texture that is inherent in the human iris."}
{"text": "1. Field of Invention\nThe present invention relates to the computer field, more particularly, to an interrupt processing method and an interrupt processing system.\n2. Description of Prior Art\nCurrently, front-end and back-end drivers of virtual machines are common architectures for virtual machine shared devices. Taking a network card as an example, in most applications, a plurality of virtual machines need to share a network. Different from a case of a graphic card, the virtual machines, even in the background, use the network. Thus, a general way is to share the network card by way of the front-end and back-end drivers.\nIn comparison to a Native system (a local system running on a local machine) driver, a front-end or back-end (FE/BE) of a virtual machine mainly differs in that a process of calling a VMCall to exchange data through a shared data region of a virtual machine manager (VMM) is increased. However, the VMCall highly consumes running time of a CPU (one VMCall needs to consume more than two thousand clock cycles, whereas a common instruction only needs one or several clock cycles), since it involves a CPU running level and context switching, switching of registers and memory and the like, and thereby it has a great influence on the performance of the whole system.\nTherefore, in addition to traditional means for optimizing device drivers (such as, changing an interrupt into a polling, changing a linked list into an array/space for time, code logical optimization, etc.), the main means to optimize the front-end and back-end of the virtual machine is to decrease the calling of VMCall.\nAs shown in FIG. 1, a present solution 1 is illustrated. According to this solution, in a current virtual machine system, a Native device generates an interrupt after receiving data. Then, the virtual machine intercepts the interrupt and sets a flag bit in a virtual CPU, and subsequently waits for a GOS (Guest OS, referring to a client system running on the virtual machine) to get into the virtual machine to check the status of the interrupt flag bit, so as to write a corresponding value into an interrupt-related register of the virtual CPU. The virtual CPU generates a simulated hardware interrupt, and after receiving the interrupt, the GOS will call the interrupt processing function according to a query of an IDT table.\nA flowchart of a communication between the front-end driver and the back-end driver is shown in FIG. 2, which comprises:\nstep S202, wherein the virtual machine receives data from applications, and at the same time probably carries out necessary data check or extraction (optional);\nstep S204, wherein a VMCall instruction is called, the VMM is switched to, and the data are copied into a shared region;\nstep S206, wherein the VMM instructs the back-end to read data from the shared region by way of interrupt injecting or interface calling; and\nstep S208, wherein the back-end driver obtains data from the shared region, and sends the data through a real device driver.\nA disadvantage of this solution lies in that: this process has to be experienced once an interrupt arrives, and the GOS performance will significantly decrease when a great deal of interrupts have been generated.\nAccording to a present solution 2, there are few devices with a function to close their interrupts, and drivers thereof may continuously read data by closing the interrupts, and then open the interrupts when there is no available data.\nA disadvantage of this solution lies in that: devices require a function to close their interrupts, but currently there are few devices with such function."}
{"text": "The present invention relates to compounds which neutralize axillary malodor.\nVarious approaches have been taken to the problem of axillary malodor. One approach has been the use of deodorants which may contain germicides and/or a fragrance. Germicides inhibit the reproduction of the bacteria which are believed to contribute to the production of axillary malodor. Fragrances are used to mask any odor produced. Another approach to the prevention of malodor is the use of antiperspirants. Antiperspirants inhibit bacteria and reduce the amount of sweat production thereby limiting the formation of substrate which gives rise to axillary malodor.\nCertain metals and metal salts have been suggested for use in the control of odors. For example, Japanese Disclosure No. 83-222011 describes a cream for removing odor from the axilla which includes copper powder and perfume in a base of cosmetic cream. The disclosure postulates that the copper powder acts on the secretion of the odor, thereby suppressing generation of the malodor.\nAnother description of the use of metal salts against odors appears in French Patent No. 1,394,875. This French patent describes the use of water-soluble iron and copper compounds to reduce or eliminate body odors. Various salts and complexes of cupric copper and ferrous iron are mentioned, including sulfates, chlorides, acetates, gluconates, citrates, tartrates, salts of ethylenediamine tetraacetic acid sodium or potassium, and iron and copper chlorophyllins. Various iron and copper phosphates are also mentioned, especially for use as food additives. The most effective combinations are indicated to be those in which both metals are present in the ionizable and nonionizable forms. The only testing reported in the French patent was odor reduction of sewage sludge, and odor reduction in mice by measurement of the odor of the entire mouse, fecal odor and urine odor.\nGerman Patent No. 1,083,503 describes the deodorizing action of metal complexes of 1,3-diketones. Listed metals include copper, nickel, cobalt, calcium, zirconium, zinc, tin, aluminum, cadmium, cerium, beryllium, magnesium, and mercury. Their use in various ways, including topical application to the axilla, is discussed. These compounds are theorized to function by interfering with the metabolism of the odor-producing microorganisms.\nBritish Patent Application No. 1,581,586 describes a sanitary foot wear article which includes a composition comprising copper, silver, or a copper-silver alloy powder dispersed in and held by water-insoluble resin binder. The British patent application indicates a belief that the metal reacts with substances secreted by the foot to produce metal salts which act as astringents and also act to prevent the growth of microorganisms.\nDespite the above, attempts have continued to find compositions which are effective against axillary malodor even after it has been formed."}
{"text": "The invention relates to ship propulsion systems of the type incorporating air actuated clutches, and more specifically to an improved control system which provides an automatic engine speed advance during the engagement cycle for the pneumatically operated clutches.\nA common form of marine propulsion system employs ahead and astern air actuated clutches for connecting the prime mover to a reversing reduction gear unit for each propeller. The air actuated clutches are engaged by the inflation of an inflatable rubber and fabric air gland bonded to an outer steel rim. Friction lining on the inner surface of the gland engages a cylindrical clutch drum when the gland is inflated. When the gland is fully deflated there is no clutch engagement, and when the gland is fully inflated there is complete clutch engagement. Between these two extremes the degree of clutch engagement corresponds to the amount of inflation of the gland.\nIn certain propulsion systems a controlled slip of the clutch is provided by controlling the degree of clutch engagement. Such controlled slip permits a lower propeller shaft speed than would be otherwise available at engine idle with full clutch engagement. This is particularly advantageous for maneuvering the ship when docking or traveling in a congested area. An example of such a controlled slip system is found in my earlier U.S. Pat. No. 3,727,737 issued Apr. 17, 1973, for \"Pressure Modulating System for Reversing Clutches and Throttle Control\".\nIn the system of my earlier patent, I provided a pneumatic clutch control assembly for a ship's propulsion system that was sequentially operated to regulate the inflation of ahead and astern air inflatable clutches and to also control the prime mover speed. The control assembly was actuated by a single throttle lever located on a pilot house control stand. Movement of the lever in one direction provided forward rotation of a propeller at a speed which increased with handle travel away from neutral. Movement of the handle in the opposite direction provided astern rotation of the propeller with speed increasing as the handle was moved farther from neutral. The center position provided a neutral setting in which the engine was disconnected from the propeller and no power was transmitted, although the engine continued to idle.\nThe single lever control of both direction and speed was accomplished in the following manner: As the lever was pivoted in either direction from the neutral, air was supplied to a selector valve which selected one or the other of the ahead and astern clutches. Thereafter, and up to a first control pressure, air pressure proportional to the position of the lever away from neutral fed through a first valve to the clutch and thereby began inflating the selected clutch. During this time the engine would remain at idle speed. After a first control pressure was reached, the first valve was piloted and it connected a second path for air to the clutch. This second path had provision for an initial programmed rate of feed of air to the clutch through a choke valve so as to softly inflate the clutch. Upon reaching a second higher control pressure, full supply air pressure was connected to the clutch. After the first control pressure was reached, the continued inflation of the clutches was not dependent upon the position of the throttle lever.\nWhen the air pressure within the clutch rose to a predetermined level, the control of my earlier patent piloted a governor valve which, in effect, connected the throttle lever control to the speed governor of the engine so that the pressure supplied to the governor directly corresponded to the position of the throttle lever and the speed could be controlled by movement of the throttle lever. The throttle lever setting determined only the final operating speed and direction and all intermediate steps of clutch engagement and inflation and engine governor speed were handled automatically by the control system.\nCertain engines used for marine propulsion will stall at idle speed when the clutch is engaged. This problem is solved by increasing the speed of the engine during the period of time in which the clutch is engaged. The increased engine speed increases the engine torque capacity and the inertia of the engine while the clutch is being engaged and thereby avoids stalling of the engine.\nPresent control systems which provide for the automatic increase in speed of the engine during the time of clutch engagement have accomplished this by automatically accelerating the engine to a specific rpm for a specific period of time. While such control systems allow for adjustment of both the specific rpm to which the engine is increased and the period of time of the engine speed increase, the available ranges of speed and time are arbitrary and are not inherently related to any specific operating conditions for the propulsion system. Thus, if the throttle lever is set for a speed which is less than the increased speed which the control will call for during clutch engagement, the engine will be accelerated to a speed beyond that at which it will run after the clutch is engaged. The time period can also continue well beyond the clutch engaging cycle thereby resulting in a surge condition where the vessel accelerates too quickly and not under control of the vessel operator.\nI have provided by the present invention an improved control system which provides for an increase in the speed of the engine during the clutch engagement, but in which the engine speed boost timing is governed specifically by the clutch engaging cycle and the engine boost speed cannot exceed the speed that the engine is set to run at after the clutch engaging cycle is completed."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and apparatus for determining the vulnerability of a structure to earthquake.\n2. Description of the Related Art\nConventionally, as methods for determining the seismic vulnerability of a structure are used (1) a method of obtaining the natural frequency of a structure by vibrating the structure, (2) a method of evaluating the antiseismic performance of a structure using a checklist, and (3) an antiseismic diagnostic method of calculating an antiseismic index through structural computation.\nIn method (1), a method of obtaining the natural frequency of a structure by vibrating the structure, a structure is vibrated using a vibration generator or a weight to give a shock, and then the natural frequency of the structure is obtained based on thus generated vibration. This natural frequency is compared with a reference value to determine the antiseismic performance of the structure.\nIn method (2), a method of evaluating the antiseismic performance of a structure using a checklist, a structure is evaluated using a checklist containing predetermined items such as conditions of the ground underneath the structure, year of construction of the structure, and structural Specifications. The antiseismic performance of the structure is determined based on total points assigned in the case of this evaluation.\nIn method (3), an antiseismic diagnostic method of calculating the antiseismic index of a structure through structural computation, any of the following three methods are used, depending on the specific case. In the first method, the ultimate strength of each layer is calculated based on the cross-sectional areas of walls and columns. In the second method, the flexural strength of each layer is obtained based on the assumption that beams and floor slabs are rigid as well as on the basis of calculated bends of individual columns and walls and calculated ultimate shear strength of the layer, and subsequently the toughness of each layer is calculated based on the relationship between this flexural strength and the ultimate shear strength, quantity of reinforcing bars, and the like. In the third method, a simple nonlinear frame analysis is conducted while taking into consideration the strength and toughness of beams.\nThe above-mentioned methods for determining the seismic vulnerability of a structure have included the following problems.\nIn the method of obtaining the natural frequency of a structure by vibrating the structure, a large-scale apparatus is required for vibrating the structure, and the structure may be damaged when the apparatus is installed or when a shock is given to the structure. Also, this method requires a reference natural frequency as a basis of comparison, such as a theoretical natural frequency, an empirically derived natural frequency, or a natural frequency measured in the past.\nIn the method of evaluating the antiseismic performance of a structure using a checklist, the evaluation is rough and susceptible to the subjectivity of an examiner; the vibration characteristics of the structure cannot be obtained quantitatively.\nIn the antiseismic diagnostic method of calculating the antiseismic index of a structure through structural computation, investigation for obtaining data for use in structural computation and entry of the thus obtained data require considerable time and expenses. Further, the diagnosis requires high-level knowledge of structural computation, so participation of a specialist is indispensable."}
{"text": "1. Field of the Invention\nThis invention relates to an image forming apparatus such as a copying machine or a printer adopting an electrophotographic recording method, an electrostatic recording method or the like, and particularly to an image forming apparatus having air blowing means for blowing air to the underside of a recording material discharged by discharging means.\n2. Related Background Art\nThere has heretofore been such an apparatus for cooling a recording material immediately behind a pair of conveyance rollers as described in Japanese Patent Application Laid-open No. 2003-208043 in order to prevent an inconvenience from occurring in the apparatus due to the heat of a recording material heated and fixed by fixing means.\nAn apparatus of this kind is shown in FIG. 6 of the accompanying drawings. As shown in FIG. 6, a recording material S having reached a high temperature by a toner image thereon being fixed by a fixing device 300 passes between a pair of conveyance rollers 130 and 140 and is discharged onto a sheet discharge tray 150, and design is made such that this recording material S receives air from a fan 200 immediately behind the pair of conveyance rollers 130 and 140.\nThereby, it becomes possible to cool the recording material S, and it is possible to prevent the phenomenon that plastic film sheets such as OHT sheets stick to each other on the sheet discharge tray 150 due to the high temperature thereof, or prevent the temperature of the sheet discharge tray 150 from rising too much to thereby adversely affect the temperature of the interior of an image forming apparatus 500. Further, in a case where the image forming apparatus 500 can form images on the two sides of the recording material S, each one side of the recording material S can be cooled, and the temperature rise in the interior of the image forming apparatus 500 can be prevented to thereby form images of good quality on the recording material S.\nNow, if in the above-described apparatus, the shape of the pair of conveyance rollers 130 and 140 is that of through rollers straight and continuous in the longitudinal direction thereof, the waving of the recording material will not occur, but yet correspondingly, the entire recording material may sometimes be rounded on the sheet discharge tray, and in some cases the stackability of the recording material has been deteriorated.\nIn order to solve this problem, the applicant has proposed in Japanese Patent Application No. 2004-315229 air blowing means set so as to prevent the rounding of a recording material by air from a fan raising and conveying the recording material.\nIn the proposed apparatus mentioned above, when recording materials are stacked up to near the air outlet of the air blowing means, although within a range substantially free of a problem, the stackability and alignment of the recording materials are somewhat disturbed by the air, and a further improvement has been desired from the viewpoint of an apparatus which is better in the handling of the recording materials."}
{"text": "This invention relates generally to security systems, and more specifically, to a multifunctional security system for defined areas.\nSecurity systems are widely used to protect property and provide personal safety. Generally, such security systems provide the protection by generating an alarm in response to one or more events, for example, an unauthorized entry, a fire, a medical emergency or a manual alarm activation. Some security systems are remotely monitored. If such a security system generates an alarm, an alarm notification signal is transmitted to a central station. Upon receiving the alarm notification signal, security service personnel at the central station respond appropriately, for example, by contacting a party at the secured location to verify the alarm or by contacting an appropriate entity. If it is appropriate to do so, the security service personnel may, upon confirmation of the alarm, contact an emergency response agency (e.g., the police department, the fire department or an emergency medical team).\nHowever, such security systems are limited in scope to providing security to a well defined area, for example, a building or cluster of buildings. These security systems also may be utilized to provide security to an outdoor area. However, these security systems are intended to be a part of a substantially permanent installation. A need still exists for a security system that is portable and rapidly deployable which still provides the functionality of the above described security systems."}
{"text": "1. Field of the Invention\nThis invention relates in general to improved carriers for electronic devices and, more particularly, to improved carriers for holding semiconductor die in a stacked relationship.\n2. Background Art\nMany electronic devices, particularly semiconductor devices, are manufactured in wafer form and then cut apart into individual die. The good die are picked out and the balance discarded. The good die are then bonded into individual packages or mounted directly onto ceramic or plastic circuit boards and/or tapes. It is frequently the case that wafer fabrication is carried out in one manufacturing location while die bonding and other assembly operations are carried out in a geographically separated location, sometimes even by a different company. Accordingly, large number of semiconductor die must be transported over large distances between the wafer manufacturing locations and the assembly locations. Since the semiconductor die are both fragile and expensive, the development of reliable and convenient shipping and handling containers is a matter of great commercial significance. These shipping and handling containers are generally referred to as die carriers.\nTwo general types of die carriers have been used in the prior art. The first involves a plastic tray which contains a large number of shallow depressions, much like an egg crate with flat bottomed compartments. One die is placed in each compartment and the plastic tray is covered with a lid. The die are generally loose within the compartment so that some jostling of the die takes place during handling and shipment. This jostling can result in chipping, scratching, or other defects. The second die carrier scheme utilizes an adhesive coated plastic tape or other tacky surface to which the die are stuck. The sticky or tacky-tape usually forms the bottom of a shallow dish. Once the die are in place, a protective cover is provided. Compartments are not needed with the tacky-tape type of container and, in general, the die do not move relative to the container during handling and shipment.\nThese prior art methods have a number of deficiencies, for example:\n(1) the volumetric efficiency of these packing arrangements is poor. The die are separated laterally within the carrier, and even though the carriers may be stacked one upon the other the number of die per unit volume is low.\n(2) the tacky-tape cannot be reused and must be discarded after shipment.\n(3) the egg crate carriers, even though reusable, are bulky and expensive to ship back to the point of origin;\n(4) these prior art carriers are difficult and/or expensive to mate with automated assembly equipment. The lateral separation of the die in the carriers requires that either the carrier or the die pick-up arm be repositioned as each die is removed. Thus, an additional two-dimensional translational motion must be provided in the handling and assembly equipment in addition to that ordinarily required for the loading, unloading, and assembly operations themselves. A further complication with respect to automated assembly is that neither the egg crate carrier nor the sticky tape carrier provides a reliable positional reference or die-to-die spacing. As a consequence the handling and assembly equipment must be able to accommodate a substantial variation in the individual die positions during loading, unloading, and assembly. This increases the cost of such equipment. Thus, a need continues to exist for improved handling and shipping containers for electronic devices, particularly semiconductor die carriers, which overcome these and other deficiencies of the prior art.\nAccordingly, it is an object of the present invention to provide an improved handling and shipping container for electronic devices, especially an improved carrier for semiconductor devices, which is rugged, easily shipped, and inexpensive to manufacture and use.\nIt is a further object of the present invention to provide an improved carrier for electronic devices, which holds the devices firmly but gently during shipment to avoid jostling and which is completely reuseable.\nIt is an additional object of the present invention to provide an improved carrier for electronic devices which has a large and variable capacity, which is compact, and which has a high volumetric efficiency.\nIt is a further objective of the present invention to provide an improved carrier for electronic devices which is easy to load and unload, and which is particularly adapted for use with automated loading and unloading equipment, and with automated assembly equipment."}
{"text": "1. Field of Invention\nThis invention relates to clutches for motor vehicles and, in particular, to an improved design and application of frictional linings for enhanced performance of automotive clutches.\n2. Brief Statement of the Prior Art\nThe typical automotive clutch has a clutch disc fixedly secured to the drive shaft and supported between the flywheel and a pressure plate. The clutch cover, which is bolted onto the flywheel includes a conical (Belleville) spring that is actuated by a clutch throw out mechanism and hydraulics to release the Belleville spring force against the pressure plate, clamping the clutch disc between the pressure plate and flywheel. The clutch disc supports frictional linings which usually have a metallic base ring on the theory that the high heat conductivity of a metal base ring would enhance heat transfer and thus reduce the tendency of the clutch linings to overheat.\nThe design of the present day automotive clutch is the result of a compromise of complex and competing design factors. Because of limitations in size and accessibility, and the requirement for a moderate pedal pressure, the maximum spring load which can be used with most clutches is from about 800 to about 2,000 pounds. The demand for increased performance of automobiles has increased the torque which is transmitted through the clutch and, accordingly, manufacturers provide clutches with clutch discs having diameters as large as practical to achieve moment arms adequate to handle the large torques transmitted by the clutches.\nMost of the car manufacturers have elected to line the entire or substantially the entire face of the clutch disc with frictional linings, apparently on the theory of the more the better. This has resulted in the unit pressures applied to the clutch disc facings being approximately 25-30 psi.\nUntil the concern over the potential hazards to the public, asbestos-containing facings were the most commonly used facings. The phase out of asbestos facings has spurred development of new materials including in particular organic composite facings. These linings frequently have optimum unit pressures from about 35 to about 65 psi, or greater, which exceed the unit pressures of current clutch designs. This has caused a less than optimum performance of clutches, often observed as a chattering of the clutch, and high wear and maintenance of the clutch disc and its frictional facings. Further, the organic composite linings are particularly sensitive to high temperatures, loosing much of their performance when their temperatures exceed about 450 degrees F.\nIt is possible to improve clutch performance by increasing the size and strength of the Belleville spring, thus applying a greater load to the clutch assembly. This approach, however, is not practical for a number of reasons. The clutches used in vehicles currently on the market have hundreds of differently sized and shaped Belleville springs and a conversion of even a small portion of the vehicles on the road and presently being manufactured would require many hundreds of thousands of dollars in tooling costs. Additionally, the higher loads imposed on the clutch assembly by the increased Belleville spring force would lead to higher thrust and would require substantially complete redesigning of the clutch throw out mechanism and hydraulics. Accordingly, the approach of increasing the load applied to a clutch is not a practical solution for retrofitting the millions of vehicles currently on the road or being marketed.\nCar manufacturers also seek higher and higher performance, leading to a new type of car, the Street Legal Performance (SLP) car which has performance characteristics which rival those of the dragsters and race cars of recent years past. The traditional approach by car manufacturers has been to increase horsepower and torque dramatically and to retool the drive train including the clutch to handle the greater power and torque. Usually clutches have been upgraded by increasing the clutch diameter and/or clamp load. This approach is expensive and usually requires retooling of related mechanical equipment.\nThere exists, today, a need for a simple and efficient conversion to increase the performance of clutches, specifically by increasing the unit pressure applied to clutch facings to the optimum value which provides maximum performance of the particular frictional facings used in the clutch. There also exists a need for a clutch lining having a design and formed of materials which protect the facing from overheating."}
{"text": "This invention relates to audio reproduction systems and in particular, to systems utilizing conventional speakers for the creation of sound waves which are reproductions of originally recorded audio signals.\nConventional electro-acoustic transducer systems for the reproduction of audio signals, generally referred to as loudspeakers, are designed to output a close approximation of the sound pressure waves generated by the original audio signal. Generally, the assumption is made during the design that the playback and amplification systems are linear with respect to the relationship between the input and output waveforms. A further assumption is generally made that the loudspeakers themselves exist in a perfect acoustic environment. In fact, they are tested in anechoic chambers during the design and production quality assurance processes. Some manufacturers provide for frequency spectrum control at the loudspeaker to compensate for the actual listening environment's effects on the sound waves after they have left the loudspeaker. Very small adjustments are feasible this way before distortion is objectionable to most listeners.\nMost loudspeakers designers depend on the listener's application of audio frequency band equalization controls on the pre-amplifier to accommodate the large range of possible listening environments. This is less than desirable means of solving the problems of variable listening environments due to the side effects of conventional equalization (EQ) systems, leading to distortion of the signal, and the dynamic nature of many listening environments over time; as in a room that may have 1 or 25 persons in it, or re-arrangement of furniture.\nSolutions have been proposed to the environment coupling problem previously described. The most commonly employed prior method has been to utilize a carefully placed microphone, near the \"average\" listening zone of a room, a signal generator or a test signal tape or disk, and a spectrum analyzer. Typically, the test tones were played back through the loudspeakers while the received spectral distribution of the signal was compared to the desired flat response curve. The difference identified was either automatically or manually applied to control the equalization settings of the audio reproduction system. Using this test procedure, environments that, for example absorb bass, could easily be detected and the proper amount of boost to such lower frequency regions of the audio spectrum could be applied.\nMore sophisticated methods have also been employed to measure the signal in the time domain. To correct for echoing environments, the departure and return times of a brief audio impulse are measured and appropriate compensation taken. The microphone typically would have to be moved many times to locate so-called hot spots, which can then be deadened with acoustically absorbent material. Acoustically \"Dead\" rooms are generally measured in the same way, with reflective materials properly placed to compensate for absorption. This process can be tedious, and generally requires a highly skilled acoustician. Even after repeated trials, some rooms remain acoustically problematical. Combining loudspeakers, environments and people has presented a hitherto unsolved problem."}
{"text": "Remote control devices are typically supplied along with consumer appliances, such as televisions, radio tuners, digital video disk players, video cassette recorders, set-top cable television boxes, set-top satellite boxes, etc. The remote control devices control the associated consumer appliances by sending infrared signals with key codes to the consumer appliances. Each such key code corresponds to a function of the selected consumer appliance, such as power on, volume down, play, stop, select, channel advance, channel back, etc. The key codes for all of the functions of a particular consumer appliance together comprise a codeset. Codesets can differ from each other not only by the bit patterns assigned to the key codes, but also by the number of bits per key code, the width of the bit pulses, the frequency of the bit pulses, as well as by the frequency of the operational signal over which the bit pulses are modulated.\nIn order to avoid the situation where a signal containing a key code operates a consumer appliance that is not selected, different codesets are used for different appliances. Although each model of consumer appliance might not have a unique codeset, there are nevertheless thousands of codesets used to operate the various types, brands and models of consumer appliances sold in the world market today.\nSome remote control devices are sold separately from consumer appliances and are preloaded with codesets that operate numerous types, brands and models of consumer appliances. These separately sold remote control devices are sometimes called “universal” remote control devices. The potential market for a universal remote control device is larger the more codesets can be loaded onto the remote control device. There is a cost, however, of providing a remote control device with additional memory in which to store additional codesets. The cost of a remote control device that can store additional codesets, however, can be reduced by more efficiently using the existing memory on the remote control device, such as the read only memory (ROM) of a microcontroller integrated circuit. For example, codesets can be stored in ROM program memory of the Z8 microcontroller available from Zilog, Inc. of San Jose, Calif.\nA microcontroller architecture is sought that allows program memory to be expanded such that codesets can be efficiently stored in program memory and expanded program memory."}
{"text": "1. Technical Field\nThe disclosure generally relates to heat dissipation, and more particularly to a heat sink assembly mounted by a clip onto an electronic package.\n2. Description of Related Art\nA heat sink is usually placed in thermal contact with an electronic package such as a central processing unit (CPU), and transfers heat through conduction away from the electronic package to prevent overheating thereof. The heat sink is secured to the electronic package by a clip.\nA related heat sink clip, shown in US Patent Application Publication No. 2005/0144764A1, includes a pressing member, an engaging member and a handle. The pressing member includes an elongated elastic section with a bend at an end thereof. The bent end defines an engaging hole therein for engaging a jut of a base which surrounds the electronic package. A slot is defined in an opposite end of the elastic section. The engaging member at an end thereof has a joining part which is provided with a fitting hole for receiving another engaging jut of the base. The engaging member at an opposite end thereof has a protrusion which passes through the slot. The handle at an end thereof, pivotally connected with the protrusion, is provided with a cam. The handle at an opposite end thereof has a pressing part to pivot the handle relative to the engaging member and extends laterally to define an extension section.\nIn securing the heat sink to the base, the handle at the pressing part is depressed to pivot the handle with respect to the pivotal joint of the handle and the engaging member. The extension section provides a hook end to fasten to a lateral side of the base. The cam presses the pressing member such that the elastic section generates tensile deformation against the heat sink to secure the heat sink to the electronic package.\nHowever, to generate a downward pressing force against the heat sink via the heat sink clip, it is necessary to press the handle downwardly with the cam pressing the pressing member to thereby secure the engaging member to the base. Pressing the handle into position requires a relatively large force, and it is not easy to mount the engaging member onto the base.\nTherefore, there is a need for a heat sink clip, which can eliminate the limitations described."}
{"text": "Repair and stabilization of the thoracolumbar region of the spine for coronal plane and sagittal plane reduction and intervertebral spacer insertion utilizes various tools and implants. These tools provide access to the thoracolumbar region, insertion of spacers, insertion of pedicle screws and other anchors and attachment of rods to these screws and anchors. Improved tools for anchor insertion, rod reduction, reduction screw tab removal and spacer insertion are desired.\nFor example, various techniques are currently being used to fuse lumbar spine vertebrae to treat specific types of spine disorders and to alleviate pain. These include the minimally invasive procedure lateral lumbar interbody fusion (LLIF) that accesses the intervertebral disc space and fuses the lumbar spine using a side (lateral) surgical approach rather than a front (anterior) or back (posterior) approach. As a minimally invasive approach, LLIF reaches the spine through several small incisions, reducing tissue trauma, scarring, post-operative discomfort and pain medication use.\nThe LLIF approach facilitates a 90° vertical reach parallel to the disc and the posterior wall of the vertebrae. Specifically in patients with coronal deformity, discs are often approached from the concave side, which permits easier access and reduces the number of incisions. However, access to L4-L5 makes an exception due to the position of the iliac crest. In patients with a coronally oblique L4-L5 disc and a high crest, a convex side approach is used. If the mid lumbar concave approach and the side of the L4-L5 approach do not correspond, a contralateral L4-L5 approach is used. After posterior fixation, table break cannot be used to open the space between the ribs and the iliac crest. Therefore, L4-L5 cannot be reached without angled measurements.\nIn addition to the limited number of disc levels that can be accessed, patient specific anatomy can provide for only a left or right approach. Before surgery, the width of the working window to each space is checked by looking at the position of the vessels and lobar plexus in axial MRI/CT slices. Depending on the axial anatomy, one side could be considered safer for access than the over side.\nAngled inserters that are connected to intervertebral spaces having a centered attachment point allow for the instrument to be used for both a left approach and a right approach. However, if the attachment point of the intervertebral spacer is off center, then two separate inserter instruments are required based on the direction of approach. These separate inserters have a fixed alignment and spacer attachment based on the direction of approach. Therefore, a single angled inserter instrument is desired that can be used for both left and right approaches during an LLIF procedure."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of topical preparations in forms such as ointment or spray which accelerate the healing of thermal injury, ultraviolet injury or radiation induced burns to the skin of humans, other mammals and other living creatures.\n2. Description of the Prior Art\nSkin is the primary protective barrier for humans. Healthy skin is a major factor in overall health. Skin can be injured by mechanical disruption (e.g., cuts and scrapes), by internal mechanisms (e.g., infiltrative disease) or by absorption of energy (e.g., heat, ultraviolet, and ionizing radiation). This study focuses on the acceleration of the healing of skin injuries caused by the absorption of various forms of energy: heat, ultraviolet and ionizing radiation.\nHeat causes thermal skin injury. Burns are classified as first, second or third degree by the extent of the injury.\nDegree Characteristics\n                1 Affects only the outer layer of the skin. Symptoms: pain, redness, and swelling. (sunburn, thermal injury)        2 Affects both the outer and underlying layer of skin. Symptoms: pain, redness, swelling, and blistering.        3 Extends into deeper tissues.         Symptoms: White or blackened, charred skin that may be numb.        \nFirst- and second-degree burns are a common first aid problem. They cause pain, redness and swelling; progressing to blistering in many cases. The discomfort during healing is the primary morbidity of thermal injury. Extensive wounds and third degree burns require expert medical intervention. Burn injury is usually independent of skin type.\nUltraviolet light causes sunburn. Not all skin is equally resistant to ultraviolet injury. A common skin phenotype classification is based on the likelihood of the skin to burn or tan.\nTypeSkin ColorCharacteristics1WhiteAlways burns, never tan2WhiteUsually burns, tans less than average3WhiteSometimes mild burn, tans more than average4WhiteRarely burns, tans more than average5BrownRarely burns, tans profusely6BlackNever burns, deeply pigmented\nWhite skin is at greater risk of ultraviolet light injury. Ultraviolet radiation injury causes pain, redness and swelling, i.e. inflammation. The discomfort during healing is the primary morbidity of sunburn.\nSunscreens, protective clothing and limiting exposure are the traditional methods to protect skin from ultraviolet induced injury. Dietary intake of carotenoids or other photoprotective compounds can increase protection of skin against injury by ultraviolet or ionizing radiation.\nSkin exposed to the high energy photons in ionizing radiation develops burns not unlike severe sunburn. Radiation-induced erythema is a comorbidity of radiation treatment in oncology.\nTopical therapy is the first line therapy for minor burns, whether heat, ultraviolet or radiation induced. Anesthetics in the topical formulations mask the pain of burns while the skin heals itself.\nHumans have used plants and plant extracts as topical medications to treat various conditions from the dawn of remembered history. The use of mosses, muds and plants are well known in the “fable and folklore” pharmacopeia. George Washington Carver recognized the medicinal value of the tomato in the early 20th century. The classical treatment for burns is the plant or an extract of Aloe vera. A review of individual studies and meta analyses reveals few positive results or little objective evidence for any positive effect, and notably, some negative effects of Aloe vera extracts on various burn injuries. Other plant derived polyphenols and extracts have been evaluated topically with mixed results. Synthetic medications, such as topical corticosteroids are effective in acute radiation dermatitis although there may be substantial side effects.\nThe following published articles have addressed the subject matter of treatment of burns and various remedies that are known or have been considered in the prior art.    1. McMillan T J, Leatherman E, Ridley A, Shorrocks J, Tobi S E, Whiteside J R Cellular effects of long wavelength UV light (UVA) in mammalian cells. J Pharm Pharmacol 60(8): 969-76, 2008.    2. Fitzpatrick T B. The validity and practicality of sun reaction types I through VI. Arch Dermatol (1988) 124: 869-871.    3. Fitzpatrick T B. (1986) Ultraviolet-induced pigmentary changes: Benefits and hazards. Curr Probl Derm 15:25-38.    4. Stahl W, Sies H. Carotenoids and flavonoids contribute to nutritional protection against skin damage from sunlight. Mol Biotechnol. 37(1):26-30, 2007.    5. Stahl W, Heinrich U, Aust O, Tronnier H, Sies H. Lycopene-rich products and dietary photoprotection Photochem Photobiol Sci. 5(2):238-42, 2006.    6. Aust O, Stahl W, Sies H, Tronnier H, Heinrich U. Supplementation with tomato-based products increases lycopene, phytofluene, and phytoene levels in human serum and protects against UV-light-induced erythema. Int J Vitam Nutr Res. 75(1):54-60, 2005.    7. Stahl W, Heinrich U, Wiseman S, Eichler O, Sies H, Tronnier H. Dietary tomato paste protects against ultraviolet light-induced erythema in humans. J. Nutr. 131(5):1449-51, 2001.    8. Césarini J P, Michel L, Maurette J M, Adhoute H, Bejot M. Immediate effects of UV radiation on the skin: modification by an antioxidant complex containing carotenoids. Photodermatol Photoimmunol Photomed. 19(4):182-9, 2003.    9. Heinrich U, Gärtner C, Wiebusch M, Eichler O, Sies H, Tronnier H, Stahl W. Supplementation with beta-carotene or a similar amount of mixed carotenoids protects humans from UV-induced erythema. J Nutr. 133(1):98-101, 2003.    10. Greul A K, Grundmann J U, Heinrich F, Pfitzner I, Bernhardt J, Ambach A, Biesaiski H K, Goilnick H. Photoprotection of UV-irradiated human skin: an antioxidative combination of vitamins E and C, carotenoids, selenium and proanthocyanidins. Skin Pharmacol Appl Skin Physiol. 15(5):307-15, 2002.    11. Stahl W, Sies H. Carotenoids and protection against solar UV radiation. Skin Pharmacol Appl Skin Physiol. 15(5):291-6, 2002.    12. Sies H, Stahl W. Non-nutritive bioactive constituents of plants: lycopene, lutein and zeaxanthin. Int J Vitam Nutr Res. 73(2):95-100, 2003.    13. Dinkova-Kostova A T. Phytochemicals as protectors against ultraviolet radiation: versatility of effects and mechanisms. Planta Med. 74(13):1548-59, 2008.    14. Weiss J F, Landauer M R Protection against ionizing radiation by antioxidant nutrients and phytochemicais. Toxicology 189(1-2):1-20, 2003.    15. Maenthaisong R, Chaiyakunapruk N, Niruntraporn S, Kongkaew C. The efficacy of aloe vera used for burn wound healing: a systematic review. Burns 33(6):713-8, 2007.    16. Vogler B K, Ernst E. Aloe vera: a systematic review of its clinical effectiveness. Br J Gen Pract. 49(447):823-8, 1999.    17. Richardson J, Smith J E, McIntyre M, Thomas R, Pilkington K. Aloe vera for preventing radiation-induced skin reactions: a systematic literature review. Clin Oncol (R Coll Radiol). 17(6):478-84, 2005.    18. Kaufman T, Kalderon N, Ullmann Y, Berger J. Aloe vera gel hindered wound healing of experimental second-degree burns: a quantitative controlled study. J Burn Care Rehabil. 9(2): 156-9, 1998.    19. Reuter J, Jocher A, Stump J, Grossjohann B, Franke G, Schempp C M. Investigation of the anti-inflammatory potential of Aloe vera gel (97.5%) in the ultraviolet erythema test. Skin Pharmacol Physiol. 21(2):106-10, 2008.    20. Boström A, Lindman H, Swartling C, Berne B, Bergh J. Potent corticosteroid cream (mometasone furoate) significantly reduces acute radiation dermatitis: results from a double-blind, randomized study. Radiother Oncol. 59(3):257-65, 2001.    21. Mnich C D, Hoek K S, Virkki L V, Farkas A, Dudli C, Laine E, Urosevic M, Dummer R. Green tea extract reduces induction of p53 and apoptosis in UVB-irradiated human skin independent of transcriptional controls. Exp Dermatol. 18(1):69-77, 2009.    22. Lincz L F, Gupta S A, Wratten C R, Kilmurray J, Nash S, Seldon VI, O'Brien P C, Bell K J, Denham J W. Thrombin generation as a predictor of radiotherapy induced skin erythema. Radiother Oncol. 90(1):136-40, 2009.    23. Carver, G W How to Grow the Tomato and 115 Ways to Prepare It For the Table. Tuskegee Institute Press, Bulletin 36, 1936.\nThe following 13 patents and published patent applications are also relevant prior art:    1. United States Published Patent Application No. 2002/0183248 to Michael J. Oldham et al. on Dec. 5, 2002 for “Method Of Using Lectins For Prevention And Treatment Of Skin Diseases And Disorders” (hereafter the “Oldham Published Patent Application”);    2. U.S. Pat. No. 6,623,769 issued to Raluca Lorant et al. on Sep. 23, 2003 for “Administration Of Lycopene For Combating Skin/Mucous Membrane Damage” (hereafter the “Lorant Patent”);    3. U.S. Pat. No. 6,780,444 issued to Alma Reza on Aug. 24, 2004 for “Method Of Making Cosmetic, Pharmaceutical, And Dermatological Compositions And Compositions Made According To The Method” (hereafter the “Reza Patent”);    4. U.S. Pat. No. 6,953,574 issued to Constantin Vladimirovich Sobol et al. and assigned to Technology Commercialization, Inc. on Oct. 11, 2005 for “Method For Producing A Fermented Hydrolyzed Medium Containing Microorganisms” (hereafter the “Sobol Patent”);    5. United States Published Patent Application No. 2005/0266018 to Benson K. Boreyko et al. on Dec. 1, 2005 for “Nutraceutical Compositions With Mangosteen” (hereafter the “Boreyko Published Patent Application”);    6. United States Published Patent Application No. 2006/0013782 to Harish Mahalingam et al. on Jan. 19, 2006 for “Use Of Active Extracts To Improve The Appearance Of Skin. Lips, Hair And/Or Nails” (hereafter the “Mahalingam Published Patent Application”);    7. United States Published Patent Application No. 2006/0105059 to Thomas James McArthur on May 18, 2006 for “Fruit And/Or Vegetable Derived Composition” (hereafter the “'0105059 McArthur Published Patent Application”);    8. U.S. Pat. No. 7,048,943 issued to Yeckezkel Barenholz et al. and assigned to Yissum Research Development Company of the Hebrew University of Jerusalem on May 23, 2006 for “Carotenoid-Loaded Liposomes” (hereafter the “Barenholz Patent”);    9. U.S. Pat. No. 7,189,419 issued to Harish Mahalingam et al. and assigned to Avon Products, Inc. on Mar. 13, 2007 for “Use Of Active Extracts To Lighten Skin, Lips, Hair, And/Or Nails” (hereafter the “Mahalingam Patent”);    10. United States Published Patent Application No. 2007/0202197 to Thomas James McArthur on Aug. 30, 2007 for “Pawpaw And/Or Peach Derived Composition” (hereafter the “'0202197 McArthur Published Patent Application”);    11. United States Published Patent Application No. 2007/0286916 to Stig Bengmark on Dec. 13, 2007 for “Synbiotic Use” (hereafter the “Bengmark Published Patent Application”);    12. European Patent No. 1,122,796 issued to Gustave Giradiere on Aug. 7, 1968 for “Composition For Making A Coating On Human Tissue” (hereafter the “Giradiere European Patent”);    13. Patent Abstract of Japan No. 2000229828 issued to Uehara Shizuka et al. and assigned to Kose Corp on Aug. 22, 2000 for “Skin Bleaching Lotion” (hereafter the “Shizuka Patent Abstract of Japan”).\nThe Oldham Published Patent Application discloses a method of using lecithin for prevention and treatment of skin diseases and disorders. The inventor believed that lecithin could be administered in a variety of forms for the delivery to dermal surfaces, either topically or subcutaneously.\nThe Lorant Patent deals with treating cosmetic effects of aging by administering topically compositions which include lycopene. The patent teaches “A method for treating an individual having cutaneous signs of aging caused by expression of proteases in the extracellular matrix comprising administering to said individual an amount of lycopene effective to substantially inhibit the expression of said proteases in the extracellular matrix of said individual.”\nThe Reza Patent discloses the concept of a skin treatment which includes a mixture of a base which consists of liquid extract from a mixture of tomatoes and apples once a spontaneous chemical reaction is initiated by the mixture of the two crushed plants is completed. Specifically, the patent discloses: “A method for manufacturing a composition comprising crushing tomatoes and apples, mixing the crushed tomatoes and apples together in a weight ratio of tomatoes to apples between 0.8 and 1.25 to produce a mixture (M), allowing the mixture (M) to stand for about twenty-four hours, and filtering the mixture (M) after the standing for about twenty-four hours to produce a filtered liquid as the composition.”\nThe Sobol Patent discloses manufacturing a skin treatment by a method of fermenting vegetable hydrolysate to create a therapeutically effective preparation. This includes various chemical combinations to treat scars, wrinkles and healing of burns. The patent discloses the following method:\n“A method for producing a fermented hydrolyzed medium containing non-pathogenic microorganisms and products of their metabolism comprising the steps of:                a) providing at least one solid food ingredient reduced to small pieces;        b) providing at least one biocompatible liquid ingredient containing at least one non-pathogenic microorganism;        c) mixing said solid food ingredient with said biocompatible liquid ingredient thereby obtaining a mixture in proportions of about 10-90% liquid to about 70-75% solid food by weight;        d) adding a sugar by mixing the sugar into the mixture at about 0.1-30% by weight; and        e) fermenting the mixture at 10-58 degrees C. until acidity reaches at least about 300 degrees Terner; whereby obtaining high acidity medium with high concentration of microorganisms and products of their metabolism.”        \nThe Boreyko Published Patent Application teaches the general concept of nutraceutical compositions including various fruits and vegetables for a variety of conditions including skin ailments.\nThe Mahalingam Published Patent Application discloses using various fruits and vegetables to improve the aesthetic appearance of the skin.\nThe McArthur Published Patent Application discloses a process for treating a fruit or a vegetable to arrive at a composition which is believed to help provide benefits to the skin. The process includes:                “A process for making a composition suitable for topical application comprising the steps of a) heating at least one fruit and/or vegetable pulp to up to a temperature in the range of about 40° C. to 100° C.; b) mixing between 1 and 40% w/w of a mild base with the heated fruit and/or vegetable pulp. There is also provided a fruit and/or vegetable derived composition suitable for topical application prepared by the above process. There is further provided a fruit and/or vegetable derived composition comprising at least one fruit and/or vegetable-derived pulp and a mild base, said composition having a pH in the range of about 7.5 to about 9.5.”        \nThe Barenholz Patent discloses treatment of skin ailments that have come from UV radiation and other types of harmful radiation. The concept is to encapsulate within liposomes substantially water-immiscible carotenoids.\nThe Mahalingam Patent discloses active extracts to lighten skin, lips, hair and/or nails. The extracts include Butea frondosa, Naringi crenulata and Stenoloma chusana. These three are added to additional extracts as set forth in Column 3 Line 54 which are Azadirachata indica, Glycyrrhiza glabra linn., Morinda citrifolia, tomato glycolipid, or any combinations thereof, can reduce melanin pigmentation of hair, skin, lips and/or nails.\nThe patent alleges:                “Lightening of hair, skin, lips and/or nails, as used in the present invention, means one or more of the following benefits is achieved. These benefits include bleaching hyper-pigmented hair, skin, lips, and/or nails; reducing age spots; evening or optimizing skin discoloration; improving the appearance of dark circles under the eyes; treating melasma, cholasma, freckles, after-burn scars, and post-injury hyperpigmentation; bleaching hair on the scalp, legs, face and other areas where bleaching and color reduction are desired; and bleaching nail stains.”        \nThe McArthur Published Patent Application has the following disclosure:                “Typically, a composition of the present invention is topically applied to an animal, preferably a human, for the treatment or prophylaxis of all epidermal disorders including psoriasis, eczema, insect bites, general epidermal irritation and redness including Rosaceae and itchiness, alopecia, circulatory disorders affecting the epidermis, sunburn, windsurf and first, second and third degree burns, healing of sores, wounds and skin infections, skin cancers including sunspots, skin melanomas, and also alleviates some gum diseases and mouth ulcers and other gum and mouth dermatological disorders. Also typically, the composition of the present invention can be topically applied to arthritic joints to alleviate pain and swelling associated with all forms of arthritis as well as general joint and muscle aches and pains. Also typically, the composition of the present invention is applied as a general pain reliever. It can also be typically applied to prevent sunburn.”        \nThe Bengmark Published Patent Application discloses a formulation which includes one antioxidant, vitamin, mineral, amino acid, peptide or protein and also includes various fruits and vegetables for use to treat burns of the skin. Specifically, the patent has the following claim:                “A method for manufacturing a formulation for the prevention or treatment of stress-induced inflammatory disorder comprising providing Pediococcus pentosaceus 16:1 (LMG P-20608), Leuconostoc mesenteroides 23-77:1 (LMG P-20607), Lactobacillus paracasei subsp paracasei F-19 (LMG P-17086), and Lactobacillus plantarum 2362 (LMG P-20606) wherein the bacterial strains are in an amount of at least 1011 CFU/ml of each of the bacteria and at least four different fibres.”        \nThe Giradiere European Patent Application discloses the treatment of irritated, blotchy and sunburned skin by using products including tomato juice. The patent application claims:                “A composition for making a coating on human tissue so as to apply thereto an active substance, the composition comprising the active substance, a first solution of a water-soluble alginate and a second solution of an appropriate metal salt which is capable of coagulating the alginate, such that when the solutions are separately applied to the tissue the alginate is coagulated to make the coating.”        \nThe Japanese Patent Application discloses a skin bleaching lotion which contains tomato juice as an active ingredient. The applications claims that the content of tomato extract is preferably 0.00005-5 wt % as a solid content.\nWhile the above prior art has extensively discussed the general concept of utilizing a tomato or portions of a tomato or products derived from a tomato to help treat the skin, none of these formulations appear to have been effective and are primarily of a theoretical nature.\nThere is a significant need for an improved topical lotion and/or spray to help relieve the pain and heal burns on the skin of both humans, mammals and other animals which do not have some of the negative side effects of prior art compositions as discussed above."}
{"text": "The present invention relates to a winch assembly.\nThe invention has been primarily developed for use in tensioning a load retaining strap for restraining freight or load on vehicles including rigid-bodied trucks, semi-trailers, flat-deck trailers, drop-deck trailers, car carriers and logging trailers and will be described with reference to these applications. However, it will be appreciated that the invention is not limited to these particular applications.\nThe use of winches in applying tension to load retaining straps in the transport industry is well known, Known winches include a winch shaft mounted for rotation in a bracket that can be attached to, for example, a semi trailer. Cylindrical extensions coupled to respective ends of the winch shaft are provided on either side of the bracket. The extensions include a radial aperture through which a drive bar or pole may be inserted to rotate the extension and thus the winch shaft.\nThe rotational axis of the winch shaft is generally substantially parallel to the longitudinal axis of the trailer and the cylindrical extension and the bar rotate about that axis. In recent times, it has become necessary to mount additional components to the trailers, for example mud guards. These components greatly restrict access to the cylindrical extension at certain locations along the length of the trailer and can make rotation by the bar or pole difficult or impossible.\nIt is an object of the present invention to substantially overcome or at least ameliorate this deficiency\nAccordingly, the present invention provides a winch assembly including:\na winch shaft adapted for rotation about a first axis; and\na substantially right angled drive transferring mechanism having a drive input coupled to a tool engagement device and drive output coupled to the winch shaft, the tool engagement device adapted for rotation about a second axis extending substantially radially from the first axis, wherein rotation of the tool engagement device causes rotation of the winch shaft.\nThe tool engagement device is preferably a nut, most preferably a nut of common size to those used in attaching road wheels to a vehicle or vehicle trailer.\nThe assembly desirably includes a mounting bracket. In one form, the bracket is adapted to be slidably received within complimentary mounting formations on a vehicle or vehicle trailer. In another form, the bracket is attached directly to the vehicle a vehicle trailer, for example by welding.\nThe drive transferring mechanism is preferably contained in a housing. In one form, the housing is fixed relative to the mounting bracket. In another form, the housing is rotatable about the first axis relative to the bracket.\nThe housing is desirably positioned adjacent one end of the winch shaft and a cylindrical extension with a radial aperture therethrough is desirably positioned adjacent the other end of the winch shaft.\nThe drive transferring apparatus preferably includes a relatively small gear coupled to the tool engagement device that meshes with a relatively large gear coupled to the winch shaft. The gears are preferably helical bevel gears.\nThe winch shaft preferably includes means for engaging a load retaining strap, most preferably in the form of a slot."}
{"text": "1. Field of the Invention\nThe present invention relates to computer system memories and, more particularly, to controlling operating voltage provided to memory devices.\n2. State of the Art\nComputer systems are typically designed to accommodate memory devices that perform within a specific band of operational parameters. For example, a computer design may accommodate specific memory devices that perform reading and writing operations at a defined speed or rate. Such an interdependent design philosophy disregards many realities of the environment of a computer system over its lifetime. For example, designing for a specific performance relationship between a microprocessor and memory devices does not allow for the independent improvements to each of the components that may, and generally does, occur. For example, microprocessor speeds may outpace memory device performance, or vice versa. In an attempt to decouple such a relationship, memory controllers have been designed to provide data brokering between the microprocessor and the memory device. Once memory controllers became ubiquitous in computer system designs, broad variations in memory device performance parameters have become commonplace.\nAdditionally, memory devices are generally tested and graded during manufacturing, with similarly performing devices integrated together into independent memory modules. As technology advances or as a computer system's memory needs change, memory modules may be upgraded or exchanged within a computer system. When memory modules are added, replaced, or exchanged with other memory modules, the memory controller adapts the timing between the memory modules and the microprocessor.\nTo date, the adaptation between the memory modules and the memory controller has been limited to modifications in timing and control parameters. However, it is known that memory technology improvements have also been made which have resulted in changes to improved or optimal operational voltages of the memory devices. Memory devices operating at a modified voltage level may exhibit an improvement in performance. Adaptation of such parameters has not been addressed by the prior art."}
{"text": "As wireless communication systems such as cellular telephones, satellite, and microwave communication systems become widely deployed and continue to attract a growing number of users, there is a pressing need to accommodate a large and variable number of communication subsystems transmitting a growing volume of data with a fixed resource such as a fixed channel bandwidth accommodating a fixed data packet size. Traditional communication system designs employing a fixed resource (e.g., a fixed data rate for each user) have become challenged to provide high, but flexible, data transmission rates in view of the rapidly growing customer base. Various standards and interoperability requirements are developed on an on-going basis for present and future communications networks. The use of standards ensures that equipment available in the marketplace operates correctly with equipment from a variety of manufacturers and service providers, and in a variety of locales so that to a user carrying the equipment from place to place and even from country to country, the use of the equipment remains convenient and the details of the operations of the network are virtually transparent to the user.\nFor example, the third Generation Partnership Project Long Term Evolution (“3GPP LTE”) is the name generally used to describe an ongoing effort across the industry to improve the Universal Mobile Telecommunications System (“UMTS”) for mobile communications. The improvements are being made to cope with continuing new requirements and the growing base of users. Goals of this broad-based project include improving communication efficiency, lowering costs, improving services, making use of new spectrum opportunities, and achieving better integration with other open standards, and backwards compatibility with some existing infrastructure that is compliant with earlier standards. The project envisions a packet-switched communications environment with support for such services as Voice over Internet Protocol (“VoIP”) and Multimedia Broadcast/Multicast Service (“MBMS”). MBMS may support services where base stations transmit to multiple user equipment (“UE”) simultaneously, such as mobile televisions or radio broadcasts, for example. The 3GPP LTE project is not itself a standard-generating effort, but will result in new recommendations for standards for the UMTS.\nThe UMTS Terrestrial Radio Access Network (“UTRAN”) includes multiple Radio Network Subsystems (“RNS”), each of which contains at least one Radio Network Controller (“RNC”). However, it should be noted that the RNC may not be present in the actual implemented systems incorporating Long Term Evolution (“LTE”) or UTRAN (“E-UTRAN”). LTE may include a centralized or decentralized entity for control information. In UTRAN operation, each RNC may be connected to multiple Node Bs, which are the UMTS counterparts to Global System for Mobile Communications (“GSM”) base stations. Generally, in this document, a base station (“BS”) is one example of a “network entity,” but many other devices that can send and receive over-the-air interface of the network are also considered a “network entities,” including other UE devices, for example. In E-UTRAN systems, the eNode B may be, or is, connected directly to the access gateway (“aGW,” sometimes referred to as the services gateway “sGW”). Each Node B may be in radio contact with multiple UE (generally, user equipment includes mobile transceivers or cellular phones, although other devices such as fixed cellular phones, mobile web browsers, laptops, personal digital assistants (“PDAs”), MP3 players, and gaming devices with transceivers may also be UE)) via the radio Uu interface. In this document, the abbreviation for user equipment (“UE”) will be synonymous with the abbreviation for mobile station (“MS”), and MS will be used primarily. Mobile stations may also be cellular phones, PDAs, MP3 players, mobile web browsers, mobile PCs and the like.\nThe wireless communication systems as described herein are applicable to, as non-limiting examples, existing wireless systems, such as 3G (“3rd Generation mobile communications”), or future systems such as 3GPP LTE compatible wireless communication systems. As a method of improving performance in such systems, use of a relay has been proposed. A network including the relay function is typically referred to as a “multiple access relay network,” or as a MARC (“multiple access relay channel”). The relay concept is that in addition to a direct transmission between a mobile station or UE and a network entity such as a base station, a relay may be used that also receives the transmission of the UE and forwards or retransmits the message, or some form of the message, to the network entity. This relay signaling would be particularly significant when a mobile station or UE signal path to a base station is less than would be desired due to obstructions such as buildings, distance, signal noise, the number of other UE in the cell or area, etc; in these situations a relay signal may be used to increase the reception at the base station and hence, the system performance. A relay may be used at any time to increase the accuracy of the reception at the network entity by providing additional signal diversity, or, redundancy. The relay is a station that listens for UE messages transmitted towards the network entity or BS, the relay then transmits a version of the signals forward towards the BS. In this manner, the BS will receive the information, or a form of the information, from the UE in the uplink direction at least twice, once from the UE itself, and once from the relay. Because signal coding operations can remove errors when multiple signals carrying the same information are received (due to signal diversity, or redundancy), when a relay is used in this manner the BS can then eliminate or reduce reception errors that might otherwise occur due to signal noise or distance from the UE.\nCoding schemes are used in communications systems. Two types of codes are of particular interest. Because the communications are performed in a relatively noisy (signal noise) environment, error correcting codes are used. Recursive systematic convolutional (“RSC”) codes are a common coding scheme. In these codes, a portion of the input is present in the output stream; also, the output is fed back into the convolution, so the code is recursive. In these codes, each m bit information symbols is input into an encoder and transformed into an n bit symbol, where n is greater than or equal to m, the code rate is m/n; and the transformation is a function of the last k information symbols, k is the constraint length of the code.\nRecent work has also focused on the turbo codes. Developed in the early 1990s, the turbo codes provide a coding scheme that performs near the theoretical Shannon limit. A turbo code requires two RSC codes and an interleaver. The decoders for turbo coders use a soft decision process, that is, a probability function for each bit is developed based on the likelihood the original bit was a “0” or a “1.” Because the decoder uses two decoders and a likelihood or estimate is made by each, an iterative process is used to change the hypotheses until the two decoders come up with the same likelihood estimates for m bits of the data payload, then the process is complete. Turbo coders are used because they provide an error correcting code that provides maximal information transfer over a communications link in the presence of data corrupting noise (for example low signal to noise ratio (“SNR”) or high error rate conditions).\nAs contemplated currently, the relay station in a relay system may simply be a UE with a different or perhaps a better signal path to the base station. The relay may be closer to the base station, or be in a path with less noise, fewer obstructions, etc. In some contemplated system arrangements, a user who allows his MS to act as a relay would receive system credit in exchange for the relay services his MS performs, which would then lower the costs of his own use of the system. In a practical device, the UE would only act as a relay if the device had plenty of charge remaining, and a good path to the BS. This feature could be disabled to conserve or extend battery life, and if the UE was busy transmitting its own data, the relay function would not be available until the device again was idle and had sufficient resources available to act as a relay.\nIn areas where signaling is known to be problematic due to noise, obstruction or distance to the nearest BS, a dedicated relay device could be used; however, in most scenarios under consideration the relay is simply an additional MS that has resources available to act as a relay. By helping other UE, each user receives a better signal and, if credit is provided, cheaper service when they use the system; and, so all of the users would benefit by having their UE act as relay stations from time to time.\nFIG. 1 depicts a simple mobile network including a relay station 13. In FIG. 1, the network entities such as base stations or node B stations 17, NB are operatively coupled together and also coupled to MME/UPE entities 18. A plurality of mobile stations MS/UE are communicating with a base station NB. In addition, a relay station RS is receiving the signals from the mobile stations MS/UE and then relaying the received signals to the base station NB. The base stations are interconnected with an X2 interface or communication link. The base stations are also connected by an S1 interface or communication link to an evolved packet core (“EPC”) including, for instance, a mobility management entity (“MME”) and a user plane entity (“UPE”), which may form an access gateway (“aGW,” a system architecture evolution gateway). The S1 interface supports a multiple entity relationship between the mobility management entities/user plane entities and the base stations and supports a functional split between the mobility management entities and the user plane entities. Note that although two UEs are shown and the relay station RS is receiving signals from two UEs, this illustration is simplified for ease of understanding and in fact many UE signals may be received by the relay and the relay will combine these received signals and relay information to the base station NB.\nThe actual implementation of such a function in the communications system continues to provide challenges. Several coding schemes have been considered for the relay. In a paper entitled “Joint Network Channel Coding for the Multiple Access Relay Channel”, Hausl et al., SECON '06, Vol. 3, September 2006, pp. 817-822, a joint network channel coding scheme, referred to hereafter as “JNCC”, is proposed. In this scheme the relay station and the UEs all use a RSC codling scheme. In another paper, entitled “A Joint Network Channel Coding Scheme for Relay Based Communications,” Heiu et al, CCECE 2007, April 2007, pp. 904-907, a similar scheme is proposed. In a paper entitled “Capacity Approaching Turbo Coding and Iterative Decoding for Relay Channels,” Zhang et al., IEEE Transactions on Communications, Vol. 53, No. 11, 2005, pp. 1895-1905; a system is proposed where the UE, if the link to the relay is less than perfect, use a turbo code at the UE and this could be used to help the relay recover the information; however, only a single UE environment is addressed.\nFIG. 2 depicts a simplified view of a JNCC system with a base station BS receiving signals from a user equipment 11, a user equipment 15, and a relay station 13. The relay station 13 also receives transmitted signals from 11 and 15. In FIG. 2, the two UEs referenced as 11 and 15 respectively transmit their information to a network entity, base station 17 with the assistance of relay station 13. In addition, the relay station 13 receives transmissions of both the UEs. Again, although only two UEs are shown transmitting signals received by the relay, many UEs would be transmitting signals received by the relay in a real system.\nIn FIG. 2, the definitions of the signal links and their corresponding signal to noise ratios (“SNR”) are as follows:                gUBi: path loss of link from a UE to the base station BS, also referred to as the “direct link”; SNRUBi is its SNR;        gURi: path loss of link from a UE to relay; SNRURi is its SNR;        gRB: path loss of link from the relay to BS; SNRRB is its SNR.        \nIn the following figures, some definitions are used to depict the coding and transmit operations:                SU,i: Systematic bits of UE i after channel coding        PU,i: Parity bits of UE i after channel coding        PUT,i: Transmitted parity bits of UE i after channel coding        PUP,i: Punctured parity bits of UE i after channel coding        SR: Systematic bits at relay after network coding        PR: Parity bits at relay after network coding        \nFIG. 3 depicts a block diagram of the encoding used in the UE of the JNCC scheme. In FIG. 3, the channel code at the UEs is a recursive systematic convolutional (RSC) code. The network code in the relay is also a RSC code with the same parameters as the RSC code used in UEs. Only the code lengths may be different.\nIn FIG. 3, symbols SU,j in data stream 35 are input to RSC channel encoder 31, and the output 37 are encoded data symbols SU,j and parity symbols PU,j These symbols are then punctured using standard code puncturing in block 33 and the output 39 are punctured, encoded RSC code symbols and corresponding parity bits.\nFIG. 4 depicts the process of a UE in the first time slot, time slot 1. In time slot 1, channel encoder (RSC) (31 in FIG. 3) encodes the systematic bits SU,i into encoded block [SU,i PU,i]. After puncturing (block 33 in FIG. 3) only SU,i and PUT,i (39 in FIG. 3) are transmitted. At the same time, the BS 17 and relay 13 listen. Note that these processes of FIG. 3 occur in both UE 11 and UE 15 in FIG. 4.\nFIG. 5 depicts the processes of the relay station 13 of FIG. 4 in the JNCC scheme in time slot 2. In time slot 2, it was assumed that the relay station RS can perfectly decode the transmitted information by UEs. At first, the relay station 13 decodes RSC encoded signals 51, 65 received from both UEs 11, 15 and recovers systematic bits SU,i referenced as 55 and 69; the decoding is done by RSC decoders 53, 67 in FIG. 5. Then the relay 13 interleaves SU,i, i=1, 2 together in interleaver 57 into a new long information block 59, which is denoted by SR. After that, the relay station encodes SR by a RSC code encoder 71 (acting as a network coding), structure of the RSC encoder is the same as in UEs and outputs encoded signals 61, which includes symbols and parity bits. Finally, after puncturing in puncturing block 73, only the newly generated parity bits PR 64 are transmitted from the relay to the network entity or BS.\nFIG. 6 depicts simply the transmission of the parity PR from the relay station 13 to the BS 17 in time slot 2.\nFIG. 7 illustrates, in a block diagram, the iterative network and channel decoder at the BS used in the JNCC scheme.\nSome definitions are needed to comprehend the blocks and operations in the FIG. 7:\nyUBi, i=1, 2: what BS received from a user equipment i in time slot 1;\nyRB: what BS received from the relay station RS in time slot 2;\n{circumflex over (x)}UBi, i=1, 2: estimated information of user equipment i after decoding at BS;\nLe−(ui)=i=1, 2: output, the extrinsic information of channel decoder;\nLe|(ui)=i=1, 2: output, the extrinsic information of network decoder.\nThe iterative decoder consists of two soft input soft output (“SISO”) channel decoders 81, 89 and one SISO network decoder 85. First, the channel decoders 81, 89 calculate extrinsic information Le−(ui), i=1, 2 of user equipment i based on received signal yUBi from UE i. A value of zero is inserted for the punctured bits PUP before decoding. The log-likelihood ratios (LLRs) Le−(ui) are interleaved and mixed in the same way as the interleaving performed in the network encoder at the relay station. The LLRs after the mixture are a priori knowledge for the network decoder.\nThe network decoder obtains additional information about the parity bits PR from the relay station by received signal yRB at the BS. The network decoder in the base station of FIG. 7 calculates extrinsic information Le|(ui) i=1, 2 about SR (which is also SU); this information is fed back to the channel decoders 81 and 89 after de-mixture and de-interleaving by functions 83, 87. After several iterations as indicated by functions 91, 93 and the switch SW1, the channel decoder at BS can combine all the available information from both the relay station and the UEs to obtain the data estimates {circumflex over (x)}UBi, i=1, 2.\nThe JNCC scheme illustrates that joint network-channel coding can improve system throughput. Information transmitted from UEs to BS and from the relay station to the BS forms a distributed turbo code, which explores the space diversity gain. In addition, network coding at the relay makes it possible for two UEs to help each other; so that once one of the two UEs is under a bad channel condition, the other UE may help its decoder to recover its original information through network decoder.\nThe known schemes of the prior art assume that the relay station can recover the UE information perfectly, but in practical communication systems this is often not the case. A continuing need thus exists for an improved coding method and apparatus to provide a robust relay coding scheme when the data estimates at the relay for signals received from the UEs is less than perfect; e.g. when there are estimation errors at the relay station, such as would typically occur in practical communications systems."}
{"text": "In the design, development and manufacture of digital circuitry, it becomes necessary to examine the characteristics of the devices that make up that circuitry, i.e., integrated circuits (ICs). Specifically, in the design and development of dynamic computer memory systems, high level standardization and characterization of dynamic MOS random access memories (RAMs), usually the heart of a dynamic computer memory system, must be maintained. An important characteristic of dynamic RAMs, in which the design engineer is interested, is its data retention time or maximum refresh rate. The maximum refresh rate is defined as the maximum time period between refresh cycles during which all locations in that dynamic RAM can maintain their validity. This characteristic is dependent upon a specified set of environmental conditions, and its response to these conditions must be carefully considered for the specific application. At the present time there is no convenient apparatus available for testing memories to determine maximum refresh rate. The present invention satisfies this need. The invention also provides a simple system for generating signals at progressively wider, controllable time intervals."}
{"text": "A microphone assembly is provided. More specifically, a ceiling mounted microphone assembly having a sensitivity pattern that is independent of the microphone's elevation angle. The microphone maximizes sensitivity in the direction of a sound source of interest, but minimizes sensitivity to sound from other directions.\nTeleconferencing systems, which can implement audio-only teleconferences or video and audio teleconferences, create meetings between two or more parties that are separately located, such as in separate rooms. The rooms may be within a same building or in different buildings, and the difference building can be located in different cities, countries, continents, etc. Thus, teleconferencing systems create meetings that would otherwise require travel of potentially large distances.\nVideo teleconferencing systems create virtual meetings by transmitting both video and audio data, and thus include one or more microphones to capture sound waves. The microphones convert sound waves generated in one video teleconferencing room into electrical impulses for transmission to another video teleconferencing room. Audio quality is therefore directly dependent on the positioning of the microphone within the room, the acoustics of the room, and to the characteristics of the microphone itself.\nFor example, a conventional microphone used to capture sound from a sound source of interest, such as a person speaking, will receive direct sound waves, reflected sound waves and reverberant sound waves from the source. Direct sound waves travel directly to the microphone without reflection, and are the sound waves intended to be captured by microphones. The level of direct sound waves is inversely proportional to the distance between the sound source of interest and the microphone receiving the sound.\nReflected sound waves do not travel directly to the microphone. Instead, they are reflected multiple times by objects in the room, or the room itself, before reaching the microphone. For example, sound waves from a sound source of interest may be reflected by walls, floors, ceilings, chairs, etc. Reflected sounds waves that propagate less than 50-80 ms (corresponding to a propagation distance of 17 to 27 meters) before reaching the microphone are known as “early reflections.” Early reflections have pressure levels approximately equal to those of direct sound waves, but are delayed in time.\nEarly reflections from the sound source of interest may positively contribute to the audio received by the microphone. However, they may also distort the audio because the time delay causes a phase difference between the early reflections and the direct sound wave that results in cancellation of part of the frequency components of the direct sound waves. This phenomenon is known as “comb filtering”, and negatively impacts quality.\nReflections that propagate for more than 50 to 80 ms (17 to 27 meters) are known as “reverberant sound”. Reverberant sound arrives at the microphone from nearly every direction because these sound waves have reflected many times within the room. Also, their pressure level is largely independent of microphone-sound-source distance. Unlike early reflections, reverberant sound always contributes negatively to audio quality by creating a “distant”, “hollow”, and/or “muffled” characteristic.\nThe level of distortion cause by reverberant sound is determined by a ratio of a level of direct sound to a level of reverberant sound. For example, if the sound source of interest is very close to the microphone the ratio of direct sound to reverberant sound is large, and distortion is small. As the sound source of interest moves away from the microphone the ratio of direct sound to reverberant sound will decrease, increasing distortion.\nA distance at which the level of the direct sound equals the level of the reverberant sound is known as the “room radius”, which can be determined for every room. As a sound source of interest moves outside of the room radius, reverberant sound dominates and distortion increases. Conversely, as the sound source moves within the room radius the direct sound dominates, and distortion decreases. Therefore, in conventional microphone systems, the sound source of interest should remain within the room radius to avoid significant audio distortion.\nDirect sound, reflected sound, and reverberant sound are not limited to the sound source of interest, and can also be present for noise sources in a video teleconferencing room. Noise sources include, for example, fan noise from ventilation systems, cooling fan noise from electronic equipment, noises from outside of the video teleconferencing room, and noises made directly on the table by people writing with pens, setting down cups, table-top computer keyboard typing, moving chairs, etc. Conventional teleconferencing system microphones receive direct, reflected and reverberant sound waves from these noise sources as well, deteriorating audio quality.\nFurther, each noise source has a different dominant component. For example, cooling fans installed on electrical equipment and noise originating from outside the video teleconferencing room primarily contribute noise in the form of reverberant sound waves. Noise generated directly on the table-top surface on which the microphone is placed contributes direct sound waves that travel parallel to the surface of the table. Some noise sources, such as ventilation systems, can also contribute multiple noise components, such as direct and reverberant sound waves.\nConventional microphones may also contribute noise in the form of an echo. An echo occurs when sound from a loudspeaker used to reproduce remote party audio is captured by the microphone and retransmitted to the remote party. Echoes also have direct, reflected and reverberant sound components, but dominance of one component over the others is determined by a loudspeaker-to-microphone distance, which is not always constant.\nEchoes are conventionally attenuated with echo cancellers, which are adaptive filters that train to a loudspeaker-microphone channel response. However, echo cancellers cannot prevent a microphone from receiving an echo. Instead, echo cancellers merely attenuate echoes already present in an audio signal. Because of their adaptive nature, echo cancellers require time to adapt to a given response, making time-invariant loudspeaker-microphone channel responses desirable. In practice, however, microphones are often repositioned during a video teleconference in order to capture audio from several different sound sources, and time-invariant loudspeaker-to-microphone channels are difficult to achieve. Thus, a conventional video teleconferencing system's echo cancellers are typically required to adapt multiple times. Moreover, echo cancellers have difficulty attenuating reverberant sound components, resulting in increased computational complexity as the level of reverberant echoes increase.\nThe above problems are exacerbated when omni directional microphones are used in video teleconferencing systems. An omni directional microphone receives audio from all directions with equal sensitivity, and thus receives direct, reflected and reverberant sounds from every sound source within the room, including noise sources. In fact, only noise sources below the conference table is attenuated because the table functions as a barrier to sound pressure waves. Though omni directional microphones are capable of capturing audio from all sound sources of interest without being repositioned, the resulting audio quality is poor because of captured noise.\nOne way to improve the quality of audio transmitted by a video teleconferencing system is to use directional microphones. Unlike omni directional microphones, a directional microphone has higher sensitivity with respect to certain directions over others, and inherently filters sound from at least some noise sources. This improves audio quality relative to an omni directional microphone, but also requires that a directional microphone be oriented to align its direction of highest sensitivity (“main axis”) toward the sound source of interest. Therefore, the directional microphone requires repositioning every time the sound source of interest changes position.\nDirectional microphones having a cardioid sensitivity pattern or a bidirectional sensitivity pattern are typically used in video teleconferencing. A microphone having a cardioid sensitivity has a directivity function given by: g(α)=1/2+1/2×cos(α), where α is the azimuth angle of a main axis with respect to horizontal. A typical cardioid microphone has a maximum sensitivity at α=0° and a minimum sensitivity at α=180°.\nA bidirectional microphone has a directivity function given by: g(α)=cos(α), where α is also the azimuth angle of a main axis with respect to horizontal. This microphone has a maximum sensitivity for α=0° and α=180°, and a minimum sensitivity when α=90° and α=270°. Because both the cardioid and bidirectional sensitivity patterns depend on the azimuth angle of the microphone, sensitivity for these microphones varies horizontally and vertically.\nAs discussed above, either a cardioid microphone or a bidirectional microphone may be used in a video teleconferencing system to improve audio quality. Placing the cardioid or bidirectional microphone on a table also improves audio quality because the table acts as a sound barrier to sound origination below the table surface, improving the direct to reverberant audio ratio.\nMicrophone sensitivity may also be improved by placing the microphone directly on the table-top surface because at this level the microphone receives direct sound waves and sound waves reflected by the table (i.e. early reflections). The direct sound waves and reflected sound waves reflected by the table, however, remain in phase and combine to form a pressure wave that is double that of the direct sound wave. This effectively increases the microphone sensitivity is by six decibels, and is commonly referred to as the “boundary principle.”\nDirectional microphones still have the drawback of requiring the sound source of interest to remain located near the main sensitivity direction of the microphone. When several people take part in the meeting, the microphone must be continually readjusted to avoid diminished audio quality. Thus, parties to the video teleconference must be aware of the sensitivity patterns of the microphone and adjust the position of the microphone accordingly. This makes directional microphones difficult to use.\nOne conventional method of reducing sensitivity to noise from the table and to orient all sound sources of interest to the “line of sight” (i.e. area of heightened sensitivity) of the microphone is to hang the microphone from the ceiling. Directional microphones, such as cardioid microphones, are often used in hanging microphone applications. However, the sensitivity pattern of hanging directional microphones is less focused than that of tabletop microphones because hanging microphones lack the shielding provided by a table surface. The missing table surface prevents hanging directional microphones from exploiting the boundary principle described above, and hanging directional microphones have relatively higher levels of self-noise compared to their tabletop counterparts. Conventional hanging directional microphones are also more susceptible to reverberant sound. Hence, conventional hanging directional microphones are only suitable for short-range use."}
{"text": "Revolvers are sometimes preferred for concealed carry applications because of their operational characteristics. Unlike semi-automatic pistols, which have a reciprocating slide that relies on delicately balanced spring pressure to chamber a fresh cartridge, revolvers have a rotating cylinder. This makes revolvers especially attractive for being carried in a pocket. Because a revolver rotates between shots, it is much less likely to get caught in clothing fabric than a semi-automatic pistol is.\nHowever, there is still a chance that the cylinder could become bound up by fabric if the revolver is fired from a pocket. Revolvers are also vulnerable to becoming locked up if the cylinder is pushed or squished.\nRevolvers have two main issues in reliability: the hammer being interfered with, and the cylinder being frozen. The former problem is known to have been addressed by partially or completely enclosing the hammer. However, the latter problem is not known to have been addressed in a serious manner.\nTherefore, a need exists for a new and improved revolver with cylinder shroud that shields the cylinder from obstruction or stoppage. In this regard, the various embodiments of the present invention substantially fulfill at least some of these needs. In this respect, the revolver with cylinder shroud according to the present invention substantially departs from the conventional concepts and designs of the prior art, and in doing so provides an apparatus primarily developed for the purpose of shielding the cylinder from obstruction or stoppage, such as by clothing when shooting from a pocket or by being grabbed by an assailant."}
{"text": "When a switch transmits a Protocol Data Unit (“PDU”) (interchangeably used herein with “packet”) across a network, the PDU may be received at a port of the same switch that sent the PDU. In such a scenario, a network loop exists. Generally, network loops are unintended, and undesired, as they negatively impact network performance."}
{"text": "1. Field of the Invention\nThe present invention relates generally to enhanced physical downlink control channel (ePDCCH) transmission, and more particularly, to a method and an apparatus for providing advanced notification of ePDCCH transmission for a user equipment (UE).\n2. Description of the Related Art\nIn the 3rd Generation Partnership Project (3GPP) Long Term Evolution (LTE) standard, a physical downlink shared channel (PDSCH) carries all user data signaling messages and a physical downlink control channel (PDCCH) carries layer one downlink control information (DCI). PDCCH is transmitted in a downlink control region and is allocated within a number of orthogonal frequency division multiplexing (OFDM) symbols at the beginning of each subframe. This allocation allows a UE to decode the PDCCH first and determine its potential assignments early, without having to read and buffer the entire subframe. Therefore, the UE can ignore the remaining portions of the subframe when the UE determines that no downlink assignment exists for that UE in the PDSCH region of the subframe. This mechanism has the potential to save power and extend battery life of the UE.\nIn order to increase the capacity of control channels to accommodate more complicated features, the concept of ePDCCH has been introduced. FIG. 1 is a diagram illustrating the incorporation of ePDCCH into a subframe. Specifically, ePDCCH regions 102 are disposed between PDSCH regions 104 in frequency, and after a PDCCH region 106 in time, within the subframe. As illustrated in FIG. 1, the ePDCCH regions 102 have a time range that begins after the PDCCH region 106 and lasts until the end of the subframe. Thus, if an ePDCCH transmission to the UE is allowed on the subframe, the UE is required to buffer the entire subframe, and then attempt to decode an ePDCCH transmission targeted to the UE. If the UE successfully decodes DCI in the PDCCH that corresponds to a downlink assignment for the UE, the UE also attempts to the decode the corresponding PDSCH transmission.\nThe battery life of the UE is negatively affected when the receiver of the UE remains active for an entire subframe and when the UE reads and buffers the entire subframe for all subframes on which an ePDCCH transmission for the UE is allowed."}
{"text": "This invention pertains generally to the field of microfluidic systems and particularly to flow sensors for such systems.\nThere are many applications where it is necessary or desirable to manipulate fluids in small volumes, including rapid bioassays, microchemical reactions, and chemical and biological sensing. See, e.g., M. Freemantle, xe2x80x9cDown Sizing Chemistry,xe2x80x9d Chem. and Eng. News, Vol. 77, No. 8, 1999, pp. 27-36. Microfluidic devices which can manipulate fluids at small scales have been formed using lithographic techniques similar to those used in microelectronic processing, often with crystalline silicon as the substrate on which the devices are formed. The microfluid flow channels that are formed in or on such substrates may have relatively small dimensions, e.g., channel widths of 1,000 xcexcm or smaller.\nFor many applications, it is necessary to measure the rate of flow in the microchannels to properly operate the microfluidic device. Not surprisingly, the design challenge of producing flow sensors which are both accurate and economical at such dimensions is formidable. One approach to the sensing of flow rates in a closed microfluidic channel is to utilize pressure sensors embedded in the walls of the channel at either end of the channel; the differential in the pressure at the two sensors is used to determine the flow rate. Such an approach has several disadvantages. It requires an especially microfabricated channel or channel cap, which can be difficult and expensive to make. In addition, it is generally difficult to accurately calibrate the pressure sensors to obtain a good reading. Further, such an approach is limited to fully enclosed channels through which the fluid flows under pressure. The typical membrane type pressure sensors formed on the walls of the microfluidic channels generally require significant additional processing steps beyond that required to produce the flow channels themselves, and the cover that extends over the substrate to close the channel typically must be micromachined to allow integration of the sensor. Another approach to flow sensing in microchannels has utilized lift forces for sensing gas flow at moderately high Reynolds numbers. See Nikalas Svedin, et al., xe2x80x9cA New Silicon Gas Flow Sensor Based on Lift-Force,xe2x80x9d JMEMS, Vol. 7, No. 3, 1988, pp. 303-308. However, such devices are not well suited for sensing fluid flow in smaller microfluidic channels and at lower Reynolds number flow.\nIn accordance with the present invention, sensing of microfluidic flow is carried out by confining and directing a fluid along a surface in a primary direction of flow past a cantilever beam which is mounted at one end of the beam to the surface. The cantilever beam has opposite beam surfaces that are oriented at an angle off parallel to the primary direction of flow of the fluid. As the fluid is directed past the beam at a rate such that the drag forces imposed by the fluid on the opposite surfaces of the beam are greater than the inertial forces of the fluid on the beam, a differential force is applied to the beam that tends to pivot the beam about its mount to the surface or bend the beam, or both. The differential drag forces apply a net force to the beam to deflect it in a direction further away from a parallel to the primary direction of flow, i.e., deflection toward an upstream direction of the flowing fluid, rather than in a direction toward the parallel to the primary direction of flow (deflection in a downstream direction), as would be the case with a beam oriented off axis in a flowing stream under flow conditions that are typically encountered in the normal macro-world where inertial forces typically dominate drag forces. However, in accordance with the present invention, the flow rates and viscosity of the fluid flowing through relatively small channels, e.g., with lateral dimensions of 500 xcexcm or less, result in drag forces dominating, making such devices well suited for microfluidic flow sensing applications. The deflection of the beam in response to the differential drag forces imposed by the flowing fluid may then be detected to determine the rate of flow of the fluid.\nMicrofluidic flow sensing apparatus in accordance with the invention may include a substrate having a channel formed therein defined by side walls and a bottom wall, and a cover on the substrate enclosing the channel and defining a top wall for the channel. The channel is formed to guide flow of a fluid therethrough in a primary direction of flow. An elongated beam having preferably parallel opposite surfaces is cantilever mounted at one end of the beam to a wall of the channel, and contained within the channel, with the surfaces of the beam oriented at an angle off of parallel to the primary direction of flow of fluid in the channel. Various means may be provided for detecting deflection of the beam in the channel in response to differential drag forces on the beam imposed by fluid flowing through the channel in the primary direction of flow. These include strain sensors formed in the wall of the channel to which the beam is mounted to sense strain in the wall caused by pivoting of the beam, photodiodes formed in the wall of the channel to which the beam is mounted at, a position beneath a portion of the beam such that the beam partially blocks light incident on the photodiodes and blocks more or less light as the beam is deflected, and an optical observation system projecting light through the cover and optically detecting deflections of the beam. Such detection means are by way of example only and any other type of detector which can detect deflection of a beam may also be utilized. The cover that is mounted to the substrate to close the channel may be formed of any of a variety of materials, including polymers, since the cover does not need to be micromachined in order to accommodate the microbeam mounted within the channel. If desired, the sensing structure can be formed on a single chip by an IC manufacturer specializing in microelectronic integration, and the cover can be formed by a microfluidic foundry specializing in fabricating plastic or elastomeric parts. Typical dimensions for the channel include a height and width of less than 1,000 xcexcm, preferably from 1 xcexcm to 500 xcexcm, a thickness of the beam less than 200 xcexcm, and a length of the beam generally at least 10 times the thickness of the beam. The substrate in which or on which the channel is formed may be various materials, including crystalline silicon, which may be machined or patterned with flow paths using well developed micromachining and patterning techniques.\nFurther objects, features and advantages of the invention will be apparent from the following detailed description when taken in conjunction with the accompanying drawings."}
{"text": "A Pb—Sn alloy plating solution has been widely used until now for plating to electronic devices for which soldering was necessary or for the formation of a solder bump electrode (bump) to a semiconductor wafer, etc. However, the Pb—Sn alloy plating solution contains toxic Ph, so that there are many problems in a wastewater treatment, environmental conservation, or a soil and groundwater pollution from wastes of the semiconductors, etc. In recent years, for the purpose of reducing such a burden on the environment, researches have been carried out as a substitute for the Pb—Sn alloy plating solution containing lead, on a plating solution using a binary alloy which utilizes tin as a first element and uses silver, bismuth, copper, indium, antimony, zinc, etc., as, a second element, or a multicomponent alloy to which a third element is further added, and among these, a Sn—Ag plating solution is now becoming a main stream. as a Ph-free plating solution.\nIn an electroplating process using the Pb—Sn alloy plating solution, however, when an anode (an anode of the electroplating process) made of Pb—Sn is used, Pb2+ and Sn2+ ions are dissolved from the anode into the plating solution, so that the balance of the components of the Pb—Sn alloy plating solution is kept substantially constant. On the other hand, in the electroplating process using, for example, a Sn—Ag plating solution, when an anode made of Sn—Ag is used, Ag is gradually precipitated on the surface of the anode to coat the surface of the anode, so that Sn2+ ion is not replenished. from the anode to the plating solution. Thus, a balance of components of the plating solution is changed, whereby it causes a problem in the electroplating process using an anode made of Sn—Ag in the Sn—Ag plating solution. Therefore, in the electroplating process using a Sn—Ag plating solution, an insoluble platinum-plated titanium plate, etc. is used as an anode.\nAlso, when an insoluble anode is used, for replenishing components of the plating solution, a method in which a metal tin is supplied by dissolving in the plating solution can be considered. According to this method, however, in the case of an alloy plating with a metal which is nobler than tin, a noble metal is precipitated on the surface of the metal tin by substituting it in the plating solution so that there is a problem that dissolution of the metal tin is suppressed. Thus, replenishment has generally been carried out by adding a tin salt solution in which essential components of the plating solution had previously been dissolved (for example, see Patent Document 1.).\nHowever, as described in the above-mentioned Patent Document 1, in the method in which Sn2+ ion is replenished by injecting a tin salt solution, etc. (hereinafter referred to as an “element solution”) where essential components of the plating solution are dissolved therein, it must prepare the element solution to be injected, and the element solution must be injected while analyzing the components of the plating solution whereby control of the plating solution is difficult and also a large cost is necessary.\nTo overcome the above-mentioned problems, it has been investigated a method in which tin (II) oxide powder having an extremely high solubility in an acid or an acidic plating solution is directly added to the plating solution whereby the tin component in the plating solution is replenished (for example, see Patent Document 2.)."}
{"text": "1. Field of the Invention\nThe present invention relates to a fluorescent material, and more particularly, to a fluorescent material used in an active dynamic liquid crystal display device and a method for manufacturing the same.\n2. Description of the Prior Art\nGenerally, in an active dynamic liquid crystal display device, a light emitted from a back light is directed to a liquid crystal layer through a polarizer. The light directed to the liquid crystal layer is transmitted to a fluorescent material or interrupted in accordance with orientations of the liquid crystal molecules. The light transmitted to the fluorescent material excites the material to display an image.\nAs described above, since the active dynamic liquid crystal display device is designed such that the light emitted from the back light directly excites the fluorescent material, there is a need for a fluorescent material with high luminance.\nIn general, SrGa.sub.2 S.sub.4 :Eu which is excited by violet light having a wavelength of about 380 to about 420 nm has been used as a green fluorescent material, and BaMgAl.sub.10 O.sub.17 :Eu has been used as a blue fluorescent material.\nHowever, the light emission luminance of these materials are not enough to apply to the active dynamic liquid crystal display device.\nTherefore, U.S. Pat. Nos. 4,249,108 and 4,733,126 have proposed a method for improving the luminous efficiency by changing composition of the fluorescent material or adding another material to the fluorescent material, and U.S. Pat. No. 5,350,971 has proposed a method for improving the luminous efficiency by adjusting a diameter of each grains of the fluorescent material.\nHowever, the fluorescent materials proposed by the above patents are not still enough in light emission luminance.\nIn addition, to increase the luminous efficiency, there has been proposed SrGa.sub.2 S.sub.4 :Eu,Tb and BaMgAl.sub.10 O.sub.17 :Eu,Tb. These materials are not also still enough in luminance.\nAs is well known, the active dynamic liquid display device requires a fluorescent material which is excited in a dark blue color range. The luminous efficiency of SrGa.sub.2 S.sub.4 :Eu and BaMgAl.sub.10 O.sub.17 :Eu is determined by transmission efficiency of light emission energy to Eu. Tb is used to transmit the light emission energy to Eu."}
{"text": "There are many and diverse types of apparatus which function to move an object in various directions. Among these apparatus are included power operated implements such as buckets, booms and bulldozer blades which are usually mounted on a vehicle and manipulated by a plurality of individual controls. The operator must give a great deal of attention to selectively moving the various levers, buttons, or the like into their proper modes of operation in order to correspondingly maneuver the various hydraulic motors or jacks which control the complex movement of the implement.\nFor example, while two separate levers can be utilized for individually operating a pair of valves in a hydraulic system and thereby directing fluid selectively to the opposite ends of a pair of hydraulic jacks, there is frequently the need to simultaneously move the jacks in the same direction or in opposite directions at the same rate. This is difficult to accomplish when the operator is also maneuvering the vehicle or manipulating a separate control element since two hands are usually required. On the other hand, if an attempt is made to mechanically connect each of the levers with both of the control valves so that the two jacks can be simultaneously extended or retracted with one lever or can be simultaneously moved in opposite directions with the other, then movement of one lever would disadvantageously move the other. In addition, it is usually desirable that both control levers be automatically returned to a neutral or jack-holding position upon manually releasing them, and yet when operating an individual lever the effort which is required should be minimal. It is thus apparent that a complex coupling problem is involved when connecting both levers to both valves.\nWhile it would appear to be desirable to utilize a single control lever to position a pair of hydraulic jacks through a suitable control valve arrangement, it is not always easy to relate the required hand movement with the desired implement movement. Furthermore, operator movement of a universally pivotable single control lever can easily lead to unequal rates of movement of the jacks to the point of requiring correctional manipulation of the lever at a crucial operating time.\nReference is herein made to U.S. Pat. No. 3,705,631 issued Dec. 12, 1972 to D. H. Seaberg, and U.S. Pat. No. 3,795,280 issued Mar. 5, 1974, the latter of which is assigned to the assignee of the present invention, which relate to control arrangements and associated hydraulic control circuits for selectively actuating separate hydraulic jacks for lifting, tilting and angling a bulldozer blade on a tractor. While such controls have been well received by the industry, it takes a special degree of dexterity and alertness to operate them. Particularly, it is difficult to correctly manipulate an electrical switch disposed on or adjacent to the single control lever. As a result, the operator will occasionally move the switch in the wrong direction or hold it engaged an excessively long period. As a consequence, one of the hydraulic jacks might well be at the end of its stroke so that the additional fluid flow being directed thereto by the inadvertent engagement of the switch must be exhausted at the relatively high pressure setting of the system relief valve. As is well known in the art, any prolonged exhausting of fluid against the system relief valve pressure is undesirable from many standpoints.\nAnother problem is that these lever control arrangements cannot always be disposed immediately adjacent to the control valves. As a result they must cooperate with various hydraulic or electrical piloting devices to remotely actuate the control valves and to thereby affect operating economies for certain member placement advantages on the vehicle."}
{"text": "1. Field of the Invention\nThe present invention relates to a fuse element used in protecting an electric circuit provided in an automobile or the like, and more particularly, to the structure of a fusible portion formed integrally with an electrical connection portion, from a thin metal plate.\nThe present application is based on Japanese Patent Application No. 2001-149914, which is incorporated herein by reference.\n2. Related Art\nA related-art automobile fuse of this type will be described by reference to that described in, e.g., Japanese Patent Application Laid-Open No. JP-A-07-65690, and shown in FIG. 5. As shown in FIG. 5, a fuse element P is formed integrally from a conductive metal plate, such as copper. A pair of connection terminals 53 are formed at the respective ends of a narrow portion 51. A thin metal layer is formed from tin on the surface of the fusible portion 51 by plating. A low-fusing-point metal chip 54 is provided in the center and at a position where the greatest amount of heat will develop when an electric current flows through the fusible portion 51. The low-fusing-point metal chip 54 is held by wrapping and caulking crimp pieces formed on the respective side edges of the fusible portion 51, thus forming a fusing portion 52. The fusible portion 51 is bent into an inverted U-shape with the fusing portion 52 taken as a center. The connection terminal portions 53 provided at the respective ends of the fusible portion 51 are connected to an electric circuit.\nIn order to cause the fusible portion 51 to satisfy a fusing characteristic with as low a resistance value as possible, there is utilized a phenomenon of the low-fusing-point metal chip 54 spreading into the fusible portion 51. A predetermined time lag is ensured by the thin metal layer formed on the surface of the fusible portion 51. However, if a comparatively low overcurrent which does not induce fusing of the fusible portion 51 within the period of the time lag flows through the fuse element repeatedly or if there has frequently arisen a phenomenon of the fusing portion 52 being instantaneously heated as a result of rushed flow of an overcurrent through the fusible portion 51, the low-fusing-point metal chip 54 induces gradual progress in diffusion alloying of the fusing portion 52. As a result, deterioration arises in the characteristic of the fusing portion 52; that is, a characteristic of shortening a time required for fusing the fusing portion 52. Hence, there arises a problem of the fuse element failing to attain predetermined durability.\nIn order to ensure a fusing characteristic and durability of a fuse, the low-fusing-point metal chip 54 requires complicated manufacturing processes and special management in relation to the purity of material, dimensional accuracy, and maintenance of fixed adhesion of the low-fusing-point metal chip 54 to the fusible portion 51, thereby resulting in low productivity and adding to manufacturing costs."}
{"text": "Two-position (on-off) HVAC actuators are typically controlled by directly applying power to a drive motor. The actuators commonly use soldered leads, printed circuit boards, conductor stampings, or a combination to make connection between the motor terminals and the output connector.\nThe polarity of the applied power determines the direction of actuator rotation. The necessary electrical connection requires two conductive paths, connecting the two connector pins to two motor terminals. The actuator housing also typically forms the shroud for the connector around the connector pins.\nThe electrical connection is usually achieved via one of the following designs approaches:\n1. Printed circuit board (PCB): a PCB is used to mount a pair of connector pins. Copper traces are screened onto the board, leading to two copper pads where motor terminals are attached (soldered). The motor is then electrically connected to these two terminals via lead wires, or directly. PA1 2. Conductor stampings: two stampings are formed into the shape of a connector pin on one end, and motor terminal on the other end. These stampings are mechanically secured into the plastic actuator housing. PA1 3. Flying leads: insulated lead wires are soldered onto the motor, and terminated with a connector. Lead wires are captured by the two halves of the actuator housing.\nIn FIG. 1, an example of a lower housing section 10 of an actuator housing is shown. A connector shroud 12 is integral with the housing section 10 and leads into the interior 14 of the housing actuator. A printed circuit board 16 is shown which fits into the lower housing section 10 such that two connector pins 18a, 18b project into the connector shroud 12.\nA free end of the connector pins forms the electrical connection between a controller or a power supply inserted into the shroud 12. The opposite ends of the connector pins are welded to the underside of the PCB 16.\nThe connector pins are electrically connected to two trifurcated terminals 20a, 20b by copper connectors screened to the underside of the PCB. The terminals 20a, 20b are also welded at a free end on the underside of the PCB.\nOptionally a capacitor can be secured to the underside of the PCB and be electrically connected between the terminals and the connector pins. The capacitor is commonly employed to reduce electromagnetic interference (EMI).\nInterior 14 of the lower housing section 10 is exposed through the connector shroud 12. This communication may allow entry of fluid to the actuator housing, and thereby potentially adversely affect the workings of the housing actuator.\nThe PCB and the other design approaches are expensive and exhibit weaknesses, ranging from allowing fluid ingress to the difficulty in adding additional components, such as a capacitor across the motor terminals."}
{"text": "The safety of the soldiers in the field is high on the priorities of weapons design. Not only are soldiers constantly exposed to enemy fire and a hostile environment during battlefield engagement, they are also exposed to the firing hazards of their own weapons and launch equipment. Often times, the launch equipment and munitions may be worn out or become unsafe due to intense use, operator error, improper maintenance or poor weather conditions, resulting in an increase of potential hazards to the soldiers. It is therefore an important goal to minimize the potential hazards to soldiers operating weapons in the field despite such real world adverse conditions.\nIn one conventional design of an artillery gun, the recoil buffering mechanism for the artillery gun comprises a breech assembly and a barrel, wherein the barrel and breech do not provide any special safety design for the operator standing behind the breech to fire the artillery gun. This configuration leaves the operator fully exposed to the dangers of explosives and hot combustion gases leaking from the gun in case the breech accidentally opens during the gun firing.\nMore specifically, when the breech is in an open position, the munition is loaded axially into the firing chamber. To perform this operation, the operator positions himself or herself at a distance from the breech. Next, the breech is closed in preparation for firing. Since the firing event is accomplished almost instantaneously, the operator remains at the distal position, behind the breech, during the entire firing operation.\nA significant pressure rise results from firing the munition. It is necessary for the breech to remain safely closed during the firing in order to impart the maximum forward momentum to the projectile, and to prevent any of the explosives and hot combustion gases from leaking past the breech to cause harm to the operator along the leakage path.\nHowever, due to wear, debris or other unforeseen factors, the breech might not be fully closed prior to, or during the firing of the weapon, resulting in leakage of explosives and hot combustion gases from the barrel to a position distal to the breech where the operator is positioned. This exposure increases the hazard of the operator and could pose substantial danger to operator safety.\nIn one embodiment, of a safety and arming mechanism for a rifled gun, the mechanism is controlled by three projectile parameters. The first and second parameters are the axial and angular accelerations of the fired projectile, which move a setback ball to arm the mechanism. The third projectile parameter, i.e., angular velocity, is utilized to lock the setback ball in the armed position. As the projectile continues its flight, it only becomes armed when a spin actuated escapement mechanism is moved to a fully armed position. In addition, a command arm signal is required to release the arrangement such that the escapement mechanism is in a condition to complete its motion to the fully armed position.\nSuch weapon arming safety design is implemented in the munition only and not in the launch equipment. It does not explicitly address safety against explosive hazards or firing hazards at the point of firing in case a catastrophic failure occurs, as in the case of firing artillery rounds, and the breech suddenly becoming loosened from the closed position, leaking explosives and hot combustion gases behind the breech.\nAnother conventional safety-and-arming device is based on micro-electromechanical system (MEMS). Two independent mechanical locks are moved out of the way to allow the arming slider to remove a barrier in the explosive train to arm a fuze or close a switch for firing. The mechanical locks respond only to valid launch or deployment conditions. In addition, the mechanism does not explicitly address safety against explosive hazards at the point of firing in case a catastrophic failure occurs.\nIn yet another conventional device, a projectile is launched with on-board linear acceleration sensors to measure at least two accelerations, and the recorded time interval between the two accelerations would need to fall within a pre-determined range in order to arm the munition for detonation. This device assures the safety of arming the munition as long as the launched projectile achieves target values in flight parameters. When this goal is not achieved, the munition in flight would not be allowed to detonate. However, this device deals with the safety to arm the projectile after becoming airborne, and not with the safety of the weapon system during the firing process to protect the weapon operators.\nIn still another embodiment, a firearm safety locking mechanism prevents accidental or unauthorized use of the weapon. The safety locking mechanism is placed and operates in the firing chamber or in the barrel of the weapon. One of the goals of this mechanism is to prevent accidental use by an under-aged operator. However, such mechanism does not address the firing hazard reduction in case the firing chamber fails to hold the hot explosive gases in place inside the weapon.\nAlthough these conventional technologies have proven to be useful, the issue of safety at the point of firing has not been addressed, and it would be desirable to present additional improvements to further reduce firing hazard. What is needed is an artillery gun equipped with a breech having a mechanism to safeguard against premature firing of munitions before the breech is fully closed. The safety mechanism should prevent explosives and hot combustion gases from the primer and the charge from quickly leaking from the firing chamber past the breech and subjecting to harm any personnel in the path of leakage. The need for such a safety mechanism has heretofore remained unsatisfied."}
{"text": "Compounds based on the ergoline ring system: ##STR1## have a suprising variety of pharmaceutical activities. For example, lysergic and isolysergic acid are D-8-carboxy-6-methyl-.DELTA..sup.9 -ergolines (9,10-didehydroergolines or 9-ergolenes.) The amides of lysergic acid have valuable and unique pharmacologic properties, and include the naturally-occurring peptide alkaloids; ergocornine, ergokryptine, ergonovine, ergocristine, ergosine, ergotamine, etc.; synthetic oxytocic alkaloids such as methergine; and the synthetic hallucinogen - lysergic acid diethylamide or LSD. Ergotamine, a 9-ergolene, with a \"peptide\" side chain, has been used in the treatment of migraine and recently, both ergocornine and 2-bromo-.alpha.-ergokryptine have been shown to be inhibitors of prolactin and of dimethylbenzanthracene (DMBA)-induced tumors in rats, according to Nagasawa and Meites, Proc. Soc. Exp'tl. Bio. Med. 135, 469 (1970) and to Heuson et al., Europ. J. Cancer, 353 (1970). (See also U.S. Pat. Nos. 3,752,888 L and 3,752,814).\nNon-peptide ergot derivatives, both naturally occurring and totally or partially synthetic, share these multiple pharmacological properties with the peptide derivatives. For example, D-6-methyl-8-cyanomethylergoline, was prepared by Semonsky and co-workers, Coll. Czech. Chem. Commun., 33, 577 (1968), and was found to be useful in preventing pregnancy in rats - Nature, 221, 666 (1969). (See also U.S. Pat. No. 3,732,231) - by interfering with the secretion of hypophysial leuteotropic hormone and the hypophysial gonadotropins or by inhibiting the secretion of prolactin. [See Seda et al., Reprod. Fert., 24, 263 (1971) and Mantle and Finn, id. 441)]. Semonsky and co-workers, Coll. Czech. Chem. Comm., 36, 220 (1971), have also prepared D-6-methyl-8-ergolinylacetamide, a compound which is stated to have anti-fertility and anti-lactating effects in rats. The 2-halo derivatives of D-6-methyl-8-cyanomethylergoline and of D-6-methyl-8-ergolinylacetamide have been prepared and found to be active prolactin inhibitors (M. J. Sweeney, J. A. Clemens, E. C. Kornfeld and G. A. Poore, 64 th Annual Meeting Amer. Assoc. Cancer Research, April, 1973 - See also U.S. Pat. No. 3,920,664).\nD-2-chloro-6-methyl-8-cyanomethylergoline (generic name, lergotrile), one of the aforementioned 2-halo derivatives, can be prepared by chlorinating D-6-methyl-8-cyanomethylergoline with such halogenating agents as N-chlorosuccinimide, N-chloroacetanilide, N-chlorophthalimide, N-chlorotetrachlorophthalimide, 1-chlorobenzotriazole, N-chloro-2,6-dichloro-4-nitroacetanilide, N-chloro-2,4,6-trichloroacetanilide, and sulfuryl chloride. Lergotrile can also be prepared by chlorinating any of the intermediate compounds used by Semonsky et al. (loc. cit.) for the preparation of D-6-methyl-8-cyanomethylergoline including methyl dihydrolysergate or D-6-methyl-8-hydroxymethylergoline, or with the novel intermediate of Example 6 of U.S. Pat. No. 3,920,664 such as, D-6-methyl-8-methylsulfonyoxylmethylergoline.\nIt is an object of this invention to provide a method of chlorinating any of the above ergoline derivatives to provide either lergotrile or a compound readily convertible thereto in higher yield than with any of the chlorinating agents heretofor provided by the art."}
{"text": "Prop shaft assemblies of the above type are often used in automotive applications to couple the transmission and/or transfer case of a vehicle to the axles in order to transfer power to the wheels.\nIn a typical prop shaft assembly, a pair of forged end yokes are provided, each having an attachment collar at one end which is press fit together with the opposite ends of a tubular drive shaft and then welded to the shaft to secure them in place. In applications where there is sufficient clearance to accommodate a large diameter drive shaft, it is preferred that the drive shaft employed be one having an inner diameter sufficiently large to accommodate the press fit insertion of the yoke collars into the ends of the shaft, as opposed to using a smaller diameter drive shaft whose outer diameter enables the ends of the shaft to be extended into the collars. One reason the large diameter shaft construction is preferred is that it is comparatively simpler and more cost effective from a manufacturing standpoint to machine the outside diameter of the yoke collars to prepare them for press fit extension and welding within the drive shaft tube, rather than having to machine the inside diameter surface of the collar to accept the shaft.\nIn some applications, however, the available space for the prop shaft assembly, and particularly the clearance for the shaft which must extend linearly between the yokes, is limited to such a degree that the preferred large diameter drive shaft construction cannot be used. A typical example of such limited clearance applications is front drive axle arrangements, where the prop shaft competes for space with the routing of the exhaust system and various other components in the vicinity of the engine compartment.\nThe solution, thus far, to such limited space requirements has been to utilize the less desirable small diameter drive shaft construction. In addition to the machining difficulties mentioned above, a small diameter drive shaft is more difficult and costly to balance. Balancing the shaft assembly involved applying weights to the end regions of the shaft to compensate for any imbalance of the assembly. The smaller diameter drive shaft offers less area on which to mount the balancing weights, as well as less net balance correction for attached weight due to the direct relationship between balance weight effectiveness and tube diameter and certain accommodations must be made for the smaller shafts since much of the standard equipment used to support the assembly and apply the weights is set up for the large diameter assemblies."}
{"text": "The disclosed subject matter relates to a gas turbine engine, and more specifically to a system for mounting a combustor transition piece to a frame of the gas turbine engine.\nA gas turbine engine includes at least one combustor, which includes a transition piece leading from a combustion zone toward a nozzle of a turbine. The transition piece is typically coupled to a frame, e.g., an aft frame, of the gas turbine engine with an integral mount (e.g., welded in place). For example, the integral mount may extend between the transition piece and the frame in a fixed orientation. Unfortunately, the integral mount may block certain areas requiring access for assembly, service, maintenance, and other operations. For example, the integral mount may block access for a welding operation. Given that the integral mount is fixed in position, the integral mount cannot be removed to provide access for these operations."}
{"text": "In general, a hearing aid operates by using a microphone to pick up sound waves in the air and convert the sound waves to electrical signals. The signals are then amplified as needed and converted back to audible sounds for the user to hear. The hearing aid microphone, however, does not always work well in conjunction with audio devices like mobile phone handsets. The acoustic connection made between the audio device and the hearing aid can be poor and can create distortions in the sound. In addition, the surrounding noise in the area of the user is often picked up by the hearing aid and interferes with the desired audio.\nMobile phones are often incompatible with many hearing aids due to the potential for radio frequency (RF) interference in the form of scattered EM fields generated by the mobile phone while in use. When a mobile phone is in communication with its network, an electromagnetic field is present around the mobile phone antenna. When the mobile phone is in use this field pulses. It is this pulsing energy that may be picked up by the hearing aid's microphone or telecoil circuitry and heard as a buzzing sound (i.e., interference) through the hearing aid.\nTo address this situation the United States Federal Communications Commission (FCC) has promulgated hearing aid compatibility (HAC) rules to minimize the interference between a mobile handset and a hearing aid. As of Sep. 16, 2005 mobile phone handset makers are to include at least 4 or 25% of their mobile phone handset offerings that meet or exceed the ANSI (American National's Standard Institute) hearing aid compatibility (HAC) standard.\nThe FCC defines HAC for mobile phones in terms of two parameters; radio-frequency (RF) emissions (M-scale) and telecoil coupling (T-scale). Mobile phones that comply with the FCC's hearing aid compatibility rule must receive a minimum rating of M3 for RF emissions and T3 for telecoil coupling. The FCC requires that nationwide carriers offer a range of phones that comply with HAC regulations beginning in September 2005 for microphone (M) mode and in September 2006 for (T) telecoil mode.\nAs part of hearing aid compatibility (HAC) assessment, near-field components of electromagnetic fields are measured. These components comprise electric and magnetic fields. Tables 1 and 2 illustrate the electric and magnetic field M ratings for the FCC HAC standards.\nTABLE 1M Ratings for HAC in standard unitsE-Field EmissionH-Field EmissionCategoryLimits (V/m)Limits A/m)M1149.6-266.10.45-0.8 M2 84.1-149.60.25-0.45M347.3-84.10.15-0.25M4<47.3<0.15\nTABLE 2M Ratings for HAC converted to dB valuesE-Field EmissionH-Field EmissionCategoryLimits (dBV/m)Limits (dBA/m)M143.5-48.5−6.9-−1.9M238.5-43.5−11.9-−6.9 M333.5-38.5−16.9-−11.9M4<33.5<−16.9\nHandsets that receive a hearing aid compatibility rating of M3 or M4 have met or surpassed the ANSI hearing aid compatibility standard as adopted by the Federal Communications Commission. An M3-rating indicates the mobile phone has satisfied the ANSI standard. An M4-rating indicates the mobile phone has exceeded the ANSI standard. The higher the M-rating the mobile phone has the lower the radio frequency (RF) emissions level and higher signal quality the mobile phone will have to the hearing aid device, assuming that the phone is still performing efficiently.\n“T” rating refers to a mobile phone's coupling ability and is intended for use with hearing aids in the telecoil mode. The higher the “T” number the more likely you will be able to use your phone with your hearing aid in telecoil mode. The FCC regulations require that cell phones be manufactured for use with hearing aids on telecoil setting by September 2006. For purposes of the present invention, however, only the M rating aspect of HAC is addressed.\nFor a phone to be HAC compliant, it must be compliant in all US frequency bands of operation for mobile phones. These include the 850-MHz and 1900-MHz bands. Moreover, the mobile phones must, in these operating bands, exhibit, at a minimum, an M3 rating meaning RF emissions from a digital wireless mobile phone are not to exceed −11.9 dB A/m for H-field and 38.6 dB V/m for E-field.\nTo the inventor's best knowledge, there has not been any mobile phone that meets the M3 rating concurrently in the 850-MHz and 1900-MHz frequency bands of operation."}
{"text": "In conventional or “wet” lithographic printing, ink receptive regions, known as image areas, are generated on a hydrophilic surface. When the surface is moistened with water and ink is applied, the hydrophilic regions retain the water and repel the ink, and the ink receptive regions accept the ink and repel the water. The ink is transferred to the surface of a material upon which the image is to be reproduced. For example, the ink can be first transferred to an intermediate blanket that in turn is used to transfer the ink to the surface of the material upon which the image is to be reproduced.\nImageable elements useful to prepare lithographic printing plates typically comprise an imageable layer applied over the hydrophilic surface of a substrate. The imageable layer includes one or more radiation-sensitive components that can be dispersed in a suitable binder. Alternatively, the radiation-sensitive component can also be the binder material. Following imaging, either the imaged regions or the non-imaged regions of the imageable layer are removed by a suitable developer, revealing the underlying hydrophilic surface of the substrate. If the imaged regions are removed, the element is considered as positive-working. Conversely, if the non-imaged regions are removed, the element is considered as negative-working. In each instance, the regions of the imageable layer (that is, the image areas) that remain are ink-receptive, and the regions of the hydrophilic surface revealed by the developing process accept water and aqueous solutions, typically a fountain solution, and repel ink.\nImaging of the imageable element with ultraviolet and/or visible radiation is typically carried out through a mask that has clear and opaque regions. Imaging takes place in the regions under the clear regions of the mask but does not occur in the regions under the opaque mask regions. If corrections are needed in the final image, a new mask must be made. This is a time-consuming process. In addition, dimensions of the mask may change slightly due to changes in temperature and humidity. Thus, the same mask, when used at different times or in different environments, may give different results and could cause registration problems.\nDirect digital imaging has obviated the need for imaging through a mask and is becoming increasingly important in the printing industry. Imageable elements for the preparation of lithographic printing plates have been developed for use with infrared lasers. Thermally imageable, multi-layer elements are described, for example, in U.S. Pat. No. 6,294,311 (Shimazu et al.), U.S. Pat. No. 6,352,812 (Shimazu et al.), U.S. Pat. No. 6,593,055 (Shimazu et al.), U.S. Pat. No. 6,352,811 (Patel et al.), U.S. Pat. No. 6,358,669 (Savariar-Hauck et al.), and U.S. Pat. No. 6,528,228 (Savariar-Hauck et al.), U.S. Patent Application Publication 2004/0067432 A1 (Kitson et al.).\nU.S. Patent Application Publication 2005/0037280 (Loccufier et al.) describes an imageable element having a single imaging layer that comprises a modified novolak resin for chemical resistance. In addition, copending and commonly assigned U.S. Ser. No. 11/204,783 (filed Aug. 16, 2005 by Ray, McCullough, Tao, and Beckley) describes the use of modified novolak resins for use in ink receptive outer layers of imageable elements."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of computer generated animation and, in particular, to a system and method for simulating object weight when generating motion of an animated character.\n2. Description of the Related Art\nUsers of game systems have come to expect increasing complexity in gameplay and increasing complexity in graphics associated with such gameplay. In many games, the visual impact on the user is as important as any other aspect of the game, creating interest in the user and aiding in explanation of the plot. Increased graphic complexity and realism additionally allows more detailed interactions within the game. The increased expectations of users of game systems for improved realism in game graphics have also increased the burden placed on animators to generate characters that are visually stimulating and to include details that were previously impossible to animate in a gameplay environment.\nIn addition, users demand that characters and objects move in a physically accurate manner and that interactions within the game are believable. In a conventional system, when a character carries an object, such as a weapon, the object usually appears to be weightless.\nAs the foregoing illustrates, there is a need in the art for an improved technique for simulating the weight of an object in an animation."}
{"text": "The present invention relates to novel polymer blends for manufacturing breathable and environmentally stable water weakenable disposable films, fibers and articles and the methods for making those blends. Particularly, the present invention relates to breathable, environmentally stable water weakenable polymer blends capable of imparting breathable and environmentally stable water weakenable films, fibers, and articles with improved stability and enhanced strength when exposed to environmental conditions of increased humidity and elevated temperature.\nDisposable products are a great convenience and provide the benefit of one time, sanitary use. Such products can be useful for applications including, but not limited to, personal care items, diapers and training pants, feminine care products, adult incontinence products, hospital garments, wound care products, hospital bed linens, surgical fabrics, medical fabrics, and the like. However, there is concern about the disposal of these products because of traditional disposal methods. For instance, incineration of these products is not desirable because of increasing concerns about air quality, and the costs and difficulties associated with separating these products from other disposed, non-incineratable articles. Dumping of these products is also undesirable due to limited landfill space and increasing land cost. Consequently, there is a need for disposable products which may be quickly and conveniently disposed of by means other than dumping or incineration.\nIt has been proposed to dispose of these products in municipal and private sewage systems. Ideally, these products could be discarded by means of water and, for example, would be water weakenable and flushable in conventional sewage systems. Articles suited for disposal in sewage systems that can be flushed down conventional toilets are termed xe2x80x9cflushable.xe2x80x9d Disposal by flushing provides the additional benefit of providing a simple, convenient and sanitary means of disposal. However, all disposable products, and especially flushable products, must have sufficient strength to perform under the conditions in which they will be used. Thus, it is desirable for disposable products to withstand the elevated temperature and increased humidity conditions encountered during use, yet lose integrity upon contact with water such as, for example, in a toilet.\nIt is also desirable that these disposable products are breathable in order to increase the level of comfort of the consumers of these products. Many disposable articles are not designed for the comfort of the user. Many of these articles use thermoplastic polymers which do not have high water vapor transmission rates and therefore do not have good breathability. In the case of many disposable products, breathability is especially desirable in order to avoid the build-up of perspiration. By increasing the breathability of the films that comprise these products, the skin wellness of the user is also increased. It is desirable for the water vapor to pass through the film and move away from the skin rather than becoming trapped against the skin where it can cause possible rashes or other skin maladies. However, it would not be desirable for the film to disintegrate when exposed to the water vapor. Therefore, it is desirable for the article to be comprised of a breathable material which has mechanical integrity when it is dry, yet it is readily water weakenable upon immersion in water. Furthermore, it is highly desirable if the film demonstrates the stability and strength to withstand the environmental conditions of elevated temperature and increased humidity that are often experienced in many end-use situations, such as, for instance, in the use of personal care products.\nPolyethylene oxide (PEO) is a commercially available thermoplastic water-soluble resin that is desirable for disposable applications. It is also desirable as a component material for flushable applications, due to its unique interaction with water and body fluids. PEO, which is represented by the following formula:\nxe2x80x94(CH2CH2O)nxe2x80x94\ncan be produced from the ring opening polymerization of the ethylene oxide, \nBecause of its water-responsive properties, PEO is desirable for flushable applications. However, commercially available PEO resins are very sensitive to high humidity and elevated temperature environments which limits its use in many disposable products.\nBecause of today\"\"s global market, it is necessary to manufacture, ship and store product components, and end-use products all over the world, encompassing a multitude of climatic conditions. PEO films dramatically lose strength and rigidity when the humidity is above about 65% Relative Humidity (RH) and at temperatures of between about 35-55xc2x0 C. because of the resulting increased moisture absorption. As a result, PEO films experience more failure and tear more easily during manufacturing and storage under these conditions of increased humidity and elevated temperatures. Especially noticeable are stress induced environmentally accelerated cracking and tearing which increases as the material ages.\nThese environmentally induced failures significantly limit the flexibility needed to manufacture, distribute, and sell disposable, flushable products. The film used for flushable products is usually folded, rolled or bent and then stored in warehouses until the end-use product is scheduled to be manufactured. The most cost effective storage space is generally not environmentally controlled and is subject to elevated temperatures and increased humidity. Such an environment accelerates the cracking and tearing of PEO films which often makes the films unsuitable for manufacture.\nShipping and storage of finished products can also cause environmentally induced tearing and cracking when the end use products are exposed to increased humidity and elevated temperature during shipment throughout the world and storage in distribution warehouses in various climates. Controlling the shipping and storage environment would significantly increase the cost of distribution.\nProducts that were subjected to increased humidity and elevated temperatures may also experience a higher incidence of failure during use causing consumer dissatisfaction. Furthermore, these problems all increase with aging of the product, limiting the length of time products could be stored before sale. These environmentally induced failures significantly limit the application of PEO, especially as a component of flushable products, because elevated temperature and increased humidity are often experienced during the manufacture, distribution and use of such products and result in performance failure.\nMany have attempted to overcome these difficulties. The current state of the art includes disposable articles consisting of a liquid impermeable, vapor permeable film consisting of a crystallizable, stretched polyolefin-based film and a rattle-reducing additive which is poly(ethylene oxide). However, these liquid impermeable, vapor permeable films require at least one nucleating agent which is described as talc or calcium carbonate. Further, stretching is required to generate porosity and hence breathability and subsequent leaching of the rattle-reducing agent is desired. These are also coated, transparent plastic articles that may consist of poly(ethylene oxide). An inorganic, protective coating is applied as a separate layer over the transparent plastic article to improve surface hardness, increase stretch resistance, and facilitate non-fogging. The inorganic, protective coating can comprise various metal oxides. However, the coating forms a separate, discrete, glass-like layer from the transparent plastic article and the resulting coating and articles are not breathable or flushable.\nMultilayered or coated recording sheets are available for electrostatic printing processes. The recording sheets consist of a base sheet with an anti-static layer, which can be made from poly(ethylene oxide). The recording sheets consist of an additional toner-receiving layer, which consists of inorganic oxides such as silicon dioxide, titanium dioxide, calcium carbonate, or the like. The poly(ethylene oxide) and inorganic oxides are contained in separate layers, the anti-static layer and the toner-receiving layer respectively. Further, the recording sheets are not breathable or flushable.\nThus, currently available polymer films are not practical for applications that are water-weakenable and breathable because they do not demonstrate improved stability and enhanced strength when exposed to environmental conditions of elevated temperatures and increased humidity. What is needed in the art, therefore, is a polymer capable of making water-weakenable and breathable films, fibers and articles that have improved stability and enhanced strength when exposed to environmental conditions of elevated temperature and increased humidity. Further, what is needed in the art is a means to efficiently and economically produce a polymer capable of forming water weakenable and breathable films, fibers and articles that have improved strength when exposed to environmental conditions of increased humidity and temperature.\nThe invention relates to novel PEO blends capable of comprising environmentally stable water weakenable and breathable films, fibers, and articles that demonstrate improved stability and enhanced strength when exposed to environmental conditions of increased humidity and elevated temperature. The invention also relates to processes and methods for manufacturing the novel blends. The novel PEO blends can be useful for manufacturing products that require disposal, including but not limited to, personal care items, diapers and training pants, feminine care products, adult incontinence products, hospital garments, wound care products, hospital bed linens, surgical fabrics, medical fabrics, and the like. Desirably, the novel PEO blends can be composed of chemically modified PEO resin and thermoplastic synthetic resins.\nThe polymer blends of the present invention can be used to produce a breathable, environmentally stable water weakenable film that does not require stretching and/or addition of filler to provide breathability. The novel blends can also provide stronger films with significantly enhanced tensile yield load, to prevent cracking and tearing when subjected to environmental stress during manufacturing, shipping and storage. Concurrently, the novel blends provide films that are water weakenable, making them desirable for disposable applications such as flushable products.\nThe invention also relates to films, fibers, and articles manufactured from the novel polymer blend compositions. These films, fibers, and articles are breathable, stable and water weakenable and demonstrate improved stability and enhanced strength when exposed to environmental conditions of increased humidity and elevated temperature. Desirably, these films, fibers, and articles can be used for disposable applications."}
{"text": "Misrepresentations by job applicants of their backgrounds put employers at risk of hiring persons lacking in training or experience that employers were counting on when the applicants were hired. Indirectly such misrepresentations, besides camouflaging the absence of training or experience, also are indicative of potentially unreliable employees. Employers also risk public embarrassment when an employee's misrepresentations are discovered, regardless of the employee's job performance. There have therefore arisen over time a variety of services providing background checking of applicants for employment."}
{"text": "1. Field of the Invention\nThe present invention relates to a method, to an apparatus and to a communications network adapted for managing electronic mail services. In particular, the present invention relates to a method, to an apparatus and to a communications network adapted for managing the sending of electronic mail messages generated by a user, for example a user having a cellular telephone. The present invention also relates, in particular, to a method, to an apparatus and to a communications network adapted for managing access to an electronic mailbox by a user, for example a user having a cellular telephone.\n2. Description of the Related Art\nElectronic mail is a system by means of which a user can exchange messages with another user by using a computer, via a communications network.\nThe use of electronic mail for exchanging messages between users of a packet network, particularly the Internet, is becoming more and more popular, owing to the great ease of use, the greater speed of message delivery (by comparison with the normal transmission of hard copy), the low cost of transmission, especially in the case of users living at great distances from each other, such as users living in different countries (by comparison with the cost of a telephone call), the possibility of attaching files of any type, such as pictures, video clips, texts, etc., to the messages to be transmitted.\nRelatively recently, it has also become possible for electronic mail messages to be transmitted and received by users of mobile networks, who can use suitable client software programs installed in some cellular telephones (known as “smartphones”, such as the model known as BlackBerry®) to manage their electronic mailboxes as if they were connected to their own home or office personal computer. In second-generation mobile networks, such as a GSM network, these cellular telephones can connect to a packet network, such as the Internet, typically via a GPRS (General Packet Radio Service) connection. Third-generation mobile networks such as UMTS (Universal Mobile Telecommunications Services) can also provide packet services. Both of these systems are therefore capable of providing an electronic mail service to users having suitable cellular telephones.\nTypically, messages sent from a sender's terminal (for example a computer or a cellular telephone) pass through one or more intermediate servers until the message is delivered to the recipient's mailbox. In an IP (Internet Protocol) environment, messages are normally transferred from the sender to the recipient by a mechanism controlled by the SMTP (Simple Mail Transfer Protocol) protocol. According to this protocol, as a result of a request to send a message made by the sender of the said message, the sender's SMTP server establishes a communications channel to a destination SMTP server, which can be the recipient's SMTP server or an intermediate SMTP server. In the latter case, the intermediate SMTP server acts as a new sending SMTP server, and the procedure is repeated until the recipient's SMTP server is reached. When the message has reached the recipient's SMTP server, it is delivered to the recipient's mailbox. The mechanism for the access to the messages present in the mailbox is typically controlled, in an IP environment, by the POP (Post Office Protocol), which typically requires the recipient to access his mailbox by means of an access key, for example a combination of a user name and password, to enable the received messages to be downloaded to his terminal.\nThe information transmitted by an electronic mail message is primarily divided into two portions: a first part, called the header, comprises the sender's address and the address of the recipient (or recipients); a second portion, known as the body, contains the actual message. Typically, the header portion of a received message also comprises a list of the SMTP servers through which the message has passed.\nThe patent application WO 99/03238 relates to a method, a system and a mail server. The method comprises the following steps: contacting an electronic mail service provider and then starting a client program to communicate with the electronic mail service. The user creates a message, determines the destination of the message and sends the message. The message is received at the server of the mail service, which requests from an exchange information relating to the user's telephone number corresponding to the user's IP address, adds the information relating to the telephone number to the header portion of the message reserved for the list of SMTP servers through which the message has passed, and sends the message towards the destination. According to the description of patent application WO 99/03238, this enables the sender of the message to be identified more precisely, particularly when the transmitted message contains “illegal” or “prohibited” information, or when it contains information which may cause the recipient's mailbox to be blocked, or when it may be convenient to contact the sender to request further information.\nThe configuration of the software program on the user terminal adapted for the use of mail services (message transmission and reception) is typically such that the user has to enter data relating to the address of the mail server or servers (for example, an SMTP server for transmission and a POP server for reception of messages), data relating to his mailbox access key (for example, a log-in and password), data relating to the mail address to be included in the sender field of the outgoing messages. These data are generally supplied to the user by the electronic mail service provider, but the user must provide by himself to the configuration of the program on the terminal. Where an electronic mail service is used by means of a cellular telephone, the configuration of the program also requires the entry of data relating to the packet connection, for example GPRS or UMTS.\nTo assist less expert users in the configuration of the program on the terminal, the electronic mail service provider generally provides instructions, which for example can be included in an Internet web page, in a paper manual or on software supports such as diskettes or CD-ROMs. In the case of access to electronic mail services by means of a cellular telephone, part of the configuration of the telephone can be carried out by pre-loading certain data into the telephone software or on to the SIM (Subscriber Identity Module) associated with the telephone, or by means of self-configuring SMS (Short Message Service) messages, or by means of procedures downloadable from the Internet, to be transferred to the cellular telephone by a PC-cellular telephone connection, such as a serial, Bluetooth or IRDA (Infra Red Digital Association) connection.\nHowever, all the above procedures for facilitating configuration only assist the users in the configuration of data which are identical for all subscribers of a service provided by a single provider (such as the data relating to the mail server addresses), while they do not permit the self-configuration of the data relating to the access key to the user's mailbox, or to the user's mail address to be included as the address of the sender of the outgoing messages: this is due to the fact that these data vary from one user to another user. Moreover, the procedures or manuals made available by the providers frequently contain technical terms which are not easily understood by non-expert users (such as the terms “server”, “username”, “login”, or the abbreviations SMTP or POP).\nThe Applicant has observed that there is a considerable proportion of the population, consisting of non-expert persons, which has difficulty in accessing electronic mail because of incorrect configuration of the electronic mail management program on the user's terminal. In some cases, these users are simply unaware of the fact that the terminal has to be configured, and, on finding that the program fails to work, they dispense with the use of the service. In other cases, the technical terminology of the configuration manual can form a real “psychological barrier”, as a result of which the user avoids using the service because he considers it too complicated.\nThe Applicant has also observed that this kind of problem particularly affects users of mobile networks. This is because these users are generally accustomed to using the cellular telephone for the ordinary telephone service. For these users, it is entirely normal to have a terminal (the cellular telephone) which is “ready for use”, without requiring more or less complex configuration procedures. The Applicant has also observed that one of the factors in the considerable success of the messaging services which have now become commonplace in mobile networks, such as SMS transmission or reception, has been the fact that the user has no need to configure anything on his mobile terminal in order to be able to use it.\nThe Applicant has tackled the problem of further facilitating access to electronic mail services, even for non-expert users. In particular, the Applicant has tackled the problem of providing an electronic mail service (comprising transmission and/or reception of messages) in which no configuration of the terminal is necessarily required to the user.\nThe Applicant has found that this problem can be solved by having the transmission and/or reception of electronic mail messages managed by at least one apparatus configured in such a way as to request, from another apparatus being part of the communications network to which a user is connected, an identifier associated uniquely with the user (for example his telephone number, or a string—so called an “alias”—associated uniquely with the user), on the basis of another identifier which the network has assigned to the user in the step of connection to the network (for example an IP address, or a host name), or which the said network is capable of associating uniquely with the user to enable them to access the network services (for example an IMSI [International Mobile Subscriber Identifier] stored in a SIM associated with the user terminal).\nTo enable an electronic mail message generated by the user to be sent, the user identifier obtained from the network can be used by the apparatus configured as stated above to automatically complete a message sender field, without the need for the user to enter it by himself among the configuration parameters. For access to a mailbox made available to the user by the provider, the apparatus configured as stated above intercepts a request for access to the mailbox by the user and uses the user identifier obtained from the network as the access key to the mailbox, without the need for the user to enter this access key by himself among the configuration parameters.\nTo be able to use an electronic mail service, a user can have, in practice as single configuration parameter, the address of the apparatus programmed as described above, in other words a fixed parameter, which is identical for all users of the service. This address can be loaded into the mail program resident in the user terminal by simple automatic configuration procedures. For example, if the electronic mail service is used from a personal computer, the address of the apparatus can be included in the configuration of the mail program by a self-configuration program made available by the service provider on a support such as a diskette or CD-ROM, or downloadable from the Internet. If the electronic mail service is used from a cellular telephone, the address of the apparatus can be included in the configuration by pre-loading on to the SIM or on to the telephone acquired by the user, or by over-the-air updating, by means of a self-configuring SMS for example."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a semiconductor memory circuit and, more particularly, to a semiconductor memory circuit having dummy cells which are connected to twisted bit lines.\n2. Description of the Prior Art\nThe current trend in the semiconductor memories art is toward higher integration density, as typified by DRAMs. The higher integration density, however, brings about a problem that the operation characteristics of a semiconductor memory circuit are adversely influenced by the coupling capacitance between bit lines which are included in the memory circuit, especially the coupling capacitance between the individual bit lines of a certain bit line pair and the adjoining bit line pairs. One approach for eliminating this problem has been reported by Yoshihara et al in a paper entitled \"Twisted Bit Line Technique for Multi-Mb DRAMs\", IEEE International Solid State Circuit Conference (1988), pp. 238-239. Specifically, in a semiconductor memory circuit having a number of bit line pairs, a number of word lines, and a number of memory cells, a twisted bit line system disclosed in the above paper is such that individual bit lines constituting a bit line pair intersect each other in places while being electrically isolated from each other. Such a configuration is successful in equalizing the coupling capacitances between the individual bit lines of a bit line pair and the adjoining bit line pairs. The paper, however, does not show how memory cells and dummy cells are arranged in the semiconductor memory circuit. A specific configuration achievable by installing dummy cells in a twisted bit line type semiconductor memory circuit on the basis of the traditional dummy cell principle would be as follows.\nFIG. 4 indicates in a schematic block diagram a part of the above-mentioned assumptive configuration, i.e., three of a number of bit line pairs and their associated elements. In the figure, the three bit line pairs are labeled OA, EA and OB. It should be noted that the structure shown in FIG. 4 is not an actual structure heretofore proposed, but it is merely an assumptive structure which would be derived from the technology disclosed in the above paper if the traditional dummy cell concept were applied thereto. In FIG. 4, the individual bit lines of each bit line pair intersect each other at places to form intersections as represented by an intersection 11, as shown and described in the paper. As used in this specification and in the appended claims, any reference to \"intersection\" or to \"intersecting\" lines or the like is not meant to denote electrical interconnection, but rather to denote a place where two lines cross one another. This is evident from the drawings which show that at these crossings, electrical contact is avoided. More specifically, assume that the entire length of bit lines is divided into four segments I, II, III and IV. Bit lines OA.sub.1 and OA.sub.2 constituting an odd bit line pair OA, bit lines OB.sub.1 and OB.sub.2 constituting an odd bit line pair OB, and so forth intersect each other in the individual odd bit line pairs at the border between the nearby segments or sections I and II and the border between the nearby sections III and IV. On the other hand, bit lines EA.sub.1 and EA.sub.2 forming an even bit line pair EA as well as bit lines forming other even bit line pairs intersect each other at the border between the sections II and III. A plurality of word lines are provided in each of the sections I to IV. In FIG. 4, two word lines a and b are shown as being laid in each of the sections I to IV for simplicity. Memory cells MC.sub.1, MC.sub.2, MC.sub.3 and so forth are individually connected to the bit lines and word lines in the junctions of the bit lines and word lines.\nIn a semiconductor memory circuit typical of which is a DRAM, dummy cells are implemented with either one of an inverted phase dummy cell method or an in-phase dummy cell method. This is to minimize the influence of the wordline-bitline coupling capacitance and to balance the individual bit lines of each bit line pair. The inverted phase dummy cell method is such that a dummy cell on a particular bit line connecting to a designated memory cell is selected by a dummy word line to which a signal whose phase is opposite to the phase of the word line is selected. On the other hand, the in-phase dummy cell method is such that a dummy cell on a particular bit line forming a pair with the other bit line which connects to a designated memory cell is selected by a dummy word line to which a signal whose phase is the same as the word line is selected.\nIn the memory circuit configuration shown in FIG. 4, the memory cells connected to the same word line are connected to those bit lines of the individual bit line pairs which correspond to each other with respect to the position in the pair. For example, the memory cells MC.sub.1, MC.sub.2 and MC.sub.3 connected to the same word line a in the section I are connected at one terminal thereof to the upper bit lines OA.sub.1, EA.sub.1 and OB.sub.1 of the bit line pairs OA, EA and OB, respectively. Then, as shown in FIG. 4, this kind of configuration would need dummy cells DC.sub.1 and DC.sub.3 and in the case of the inverted phase dummy cell method DC.sub.6 or would need dummy cells DC.sub.2, DC.sub.4 and DC.sub.5 in case of the in-phase dummy cell method having at least two of four dummy word lines A, B, C and D.\nFIG. 5 tabulates memory cells and dummy cells which are selected in the above-described semiconductor memory circuit by the in-phase dummy cell system and the anti-phase dummy cell system. Assume that the word line a existing in the section I is selected to access the memory cells MC.sub.1, MC.sub.2 and MC.sub.3 for data being stored therein. Then, the data associated with the memory cells MC.sub.1, MC.sub.2 and MC.sub.3 exist on the bit lines OA.sub.1, EA.sub.1 and OB.sub.1, respectively. It follows that in the in-phase dummy system the dummy word lines B and C are selected by a signal having the same phase as the word lines to in turn select the dummy cells DC.sub.2, DC.sub.5 and DC.sub.4, while in the antiphase dummy cell system the dummy word lines A and D are selected by a signal opposite in phase to the word lines to in turn select the dummy cells DC.sub.1, DC.sub.6 and DC.sub.3.\nIn the semiconductor memory circuit shown in FIG. 4, the region (sections I to IV) where the bit lines intersect each other and the region (indicated by a dash-and-dot line in the figure) where the dummy cells are provided are independent of each other. The circuit, therefore, cannot be integrated unless it is provided with exclusive areas for accommodating the two different kinds of regions. This obstructs efficient circuit design which is desirable for higher integration density.\nThat only a single pair of dummy cells is provided on each bit line pair as shown and described gives rise to another problem. Specifically, in each bit line pair, the time constants of the individual bit lines with respect to their associated sense amplifier are not the same as each other and, therefore, unbalanced. For example, assuming that the word line a in the section I is selected, then the dummy word lines B and C are selected in the in-phase dummy cell system and, as a result, the memory cells and the dummy cells selected are located at opposite remote ends of the bit lines. In this condition, the distance to the memory cells and the distance to the dummy cells as measured from a sense amplifier which is associated with the bit lines greatly differ from each other, resulting in the bit lines of the bit line pair differing from each other in time constant. To promote efficient sense amplification in a bit line pair, it is preferable that paired bit lines have time constants which are as close to each other as possible."}
{"text": "Conventionally, as a device for easily analyzing characteristics of bodily fluid of a man or an animal, for example, a device for electrochemically measuring blood sugar has already been put to practical use.\nA biosensor is a device of this type, and hereinafter, a biosensor and a measuring device to be combined with the biosensor will be described.\nFIG. 21 shows a long and narrow strip-shaped sensor 31 for collecting blood, which is set in a sensor insertion slot 31a of a measuring device 32. The sensor 31 is provided with a cavity (not shown) for collecting blood, in its semi-circular front end that protrudes from the measuring device 32 in a state where the sensor 31 is housed in the measuring device 32. Further, the sensor 31 is provided with, in the cavity, a reagent layer containing an enzyme, an electron carrier, and the like, and electrodes.\nThe measuring device 32 contains an electric circuit for measuring a current value according to concentration of glucose in blood, which current is generated by a reaction between the glucose in blood and the reagent layer, by applying a voltage to the electrode, and measured blood sugar is displayed on a display 33 which is placed on an upper surface of the measuring device 32.\nFor every measurement, a new sensor 31 is inserted into the sensor insertion slot on a side surface of the measuring device 32, and blood of a patient is applied onto the sensor 31 to perform measurement. The sensor 31 after measurement is discarded for hygienic reasons.\nUsually, skin of a body part, such as a fingertip, is lanced with a lancet device 34 as shown in FIG. 22(a) to extract a very small quantity of blood, and the blood is collected in the cavity of the sensor 31.\nIn the lancet device 34 whose internal structure is shown in FIG. 22(b), a force is elastically applied to a lancet 35 by a spring coil 36. When an operation button 37 is pushed, an engagement member 37a which is united with the operation button 37 is disengaged from a ring-shaped groove 35a of the lancet 35, whereby a force applied to the lancet 35 by the spring 36 is released, and a tip of a needle 35b of the lancet 35 forcibly projects from a front end of an approximately cylindrical-shaped case 38. The lancet 35 comprises the needle 35b which is made of metal, and a holder 35c for holding the needle 35b, which is made of plastic. Usually, the lancet 35 is replaced with a new one at every measurement, for hygienic reasons.\nFurther, a biosensor to be used for analyzing bodily fluid extracted from a man or an animal is usually preserved in a state where it is wrapped with an aluminum wrapper or the like, or in a state where it is contained in a plastic case. When the biosensor is used, it is taken out of the aluminum wrapper or the plastic case.\nFIG. 23(a) is a diagram illustrating a state where a biosensor 110 having a reagent layer (not shown) for analyzing bodily fluid, and an electrode (not shown) for removing an electric signal according to result of analysis, is hermetically wrapped with an aluminum wrapper 120. To be specific, a biosensor 110 which comprises an approximately rectangular-shaped plate member having a shorter side of an approximately semi-circular shape is wrapped with an aluminum wrapper 120 having a rectangular shape that is a little larger than the biosensor 110. FIG. 23(b) shows a state where a plurality of biosensors 110, which are arranged with their approximately-semi-circular-shaped ends facing upward and their surfaces being in contact with each other, are hermetically contained in a cylindrical plastic case 130 with a lid.\nWhen performing measurement, the biosensor 110 which is preserved as shown in FIG. 23(a) or 23(b) is removed, and another shorter side of the biosensor 110, which is not the approximately-semi-circular-shaped side, is inserted into a biosensor insertion slot 114a of a measuring device 114 to prepare for measurement as shown in FIG. 24.\nSince the conventional biosensor and measuring device thereof are constructed as described above, when preparing for measurement, a patient sets a new lancet 35 in the lancet device 34 as described above. Then, the patient sets a new sensor 31 in the measuring device 32, whereby preparation for measurement is completed. Thereafter, the patient operates the lancet device 34 to take blood from his/her fingertip or the like, and applies it onto a front end of the sensor 31 that is set in the measuring device 32 to perform measurement. In this way, the patient must perform replacement of the lancet and replacement of the sensor separately for every measurement, resulting in complicated operation.\nFurther, since it is necessary to monitor blood sugar a few times a day, sizes of these devices are reduced in consideration of portability. In a conventional system, however, the patient must carry the sensor 31, the measuring device 32, the lancet 35, and the lancet device 34 together, resulting in a voluminous system as a whole. Further, since the patient must manage the sensor 31 and the lancet 35 separately, operation and measurement are troublesome, resulting in poor usability.\nThe present invention is made to solve the above-mentioned problems and has for its object to provide a lancet-integrated sensor in which a sensor and a lancet are integrated to facilitate operation and management, and improve portability, as well as a measuring device to be combined with the lancet-integrated sensor.\nFurther, with respect to the conventional biosensors, as described above, each biosensor is wrapped with an aluminum wrapper, or plural biosensors are stored in a plastic case with a lid, whereby the biosensors are prevented from contamination due to moisture.\nHowever, when preparing for measurement, a patient must break the aluminum wrapper 120 to take the biosensor 110, or open the lid 130a of the plastic case 130 to take the biosensors 110 one by one. Further, the patient must insert this taken biosensor 110 toward the sensor insertion slot 114a of the measuring device 114. Therefore, preparation for measurement is troublesome, and usability is poor.\nThe present invention is made to solve the above-mentioned problems and has for its object to provide a biosensor cartridge by which a biosensor stored in a case is easily inserted into a measuring device without troublesome operation."}
{"text": "I. Field of the Present Invention\nThe present invention relates to a tube or conduit with strings or threads which is applicable to irrigation.\nII. Description of the Prior Art\nTraditional irrigation systems for plants, flower, and vegetables etc. are the tube drip irrigation system (FIGS. 1 and 2). The main problems of the traditional tube drip irrigation system are that they are not designed to solve water and fertilizer loss problems, also producers have to take care of them more frequently because the tubes are easy to drip out from pots. Seeping irrigation for garden is seeping irrigation hosepipe, (FIG. 3). The main problem of the traditional seeping irrigation hosepipe in using is that its limitation because this system can not be used to irrigate the potting plants.\nThere is a need for a new seeping irrigation system to solve the above problems. The present invention provides a solution to solve the above and other relating problems."}
{"text": "Coalbed methane (CBM) is the composite of components that may be adsorbed on coal at the naturally occurring conditions of reservoir pressure and temperature. As pressure is reduced, the CBM begins desorbing from the coal once the critical desorption pressure (CDP) is reached. CBM may consist largely of methane with smaller amounts of impurities, typically nitrogen and carbon dioxide and some minor amounts of intermediate hydrocarbons.\nThe capture and sale of CBM is a burgeoning industry both in the United States and internationally. In the CBM industry, a typical procedure for CBM recovery is often to penetrate the geologic formation with a substantially vertically drilled well and to either 1) case the hole, typically with steel casing through the coal interval followed by cementing the casing in place and perforating the interval all by methods commonly known in the petroleum industry, or 2) to case in a like manner the hole to the top of the coal and then drill through the coal, perhaps widening the hole drilled through the coal by a process known in the industry as underreaming. The former case is known as a cased completion and the latter is known as an open-hole completion. In either case, when producible water is present, typically water is pumped from the well through a tubing string to the surface in an attempt to lower the reservoir pressure, a generally necessary condition for releasing commercial quantities of CBM in most production scenarios. As reservoir pressure is lowered, a free gas phase will eventually form at the bottom of the hole and most of the free gas then will rise in the annulus between the casing and the tubing by gravitational forces, allowing the relatively buoyant gas to be produced at the surface from the annulus of the casing. The gas produced is then gathered and then typically sent to markets through pipelines.\nMany CBM wells that will ultimately produce commercial quantities of coalbed methane do not do so when first put into production. The only gas produced initially in such wells is the relatively minute, generally noncommercial, quantity of gas that is in solution in the water at bottom-hole conditions of pressure and temperature. Most of this minute quantity will come out of solution as the produced formation water moves from conditions at the bottom of the hole to the lower pressure and typically different temperature at the surface. Such coal formations that do not produce gas initially beyond the amount contained in solution in the formation water are said to be undersaturated at reservoir conditions of pressure and temperature. Other definitions for undersaturated coals include: 1) when the storage capacity of the coal, typically expressed in standard (usually 14.7 psia and 60 deg F.) cubic feet of gas per ton of coal, exceeds the actual gas content of the coal expressed in the same units at reservoir pressure, or 2) when no free gas phase exists in the cleats and fracture system at reservoir conditions.\nStorage capacity of the coal is typically determined in the laboratory from a captured sample of coal. A plot of the data is often made having the ordinate typically expressed in SCF/Ton and the abscissa being absolute pressure. This data is also often statistically fit with an equation to yield a curve, one such commonly used curve being known as the Langmuir isotherm as described in the reference of Yee et al., 1993. These “isotherms”, as the term implies, are measured at constant temperature generally corresponding to that of the reservoir from which the sample was obtained.\nUnfortunately, some of the undersaturated CBM reservoirs may never produce commercial quantities of coalbed methane. One concern, therefore, is the determination of whether or not the coals in these undersaturated CBM reservoirs contain sufficient gas to be commercial. Such information, if it could be determined expediently on a given well in an exploratory area, could prevent the drilling of a large number of wells in the specific area that may never produce economic quantities of CBM. As mentioned above, one common method of making that determination is through the process of obtaining a sample of the coal itself, perhaps by coring the coal, and subsequent detailed measurement of gas content of that sample in a laboratory or otherwise. This technique is typically expensive, and can require specialized drilling equipment and personnel. Additional expense may be incurred when the core samples are sent to commercial or private laboratories for analysis. The results of such core analyses are not immediately available, sometimes taking months of desorption time. Also, because core analysis may be too expensive for a large amount of sampling to be taken from a particular well, samples, hoped to be representative, are often selected. Consequently, there is the potential problem of the core samples not being representative of the formation even nearby the well from which the core was cut; and there is an additional problem of how representative the samples will be of the formation at some distance from the well. The CBM industry is replete with examples of how gas content can drastically change over relatively short distances. It is typically neither economically practical nor timely to have every well cored and analyzed.\nThe results from a sample of the coal itself, perhaps from the coring process, can also be very inconsistent from what is ultimately observed during production. During a coring or other sampling operation, not only are samples of coal pulled for determining gas content in the laboratory, but also a specific sample or a composite sample, possibly made up from drill cuttings, may be gathered and this sample used to determine storage capacity of the coal. This can involve tedious and expensive laboratory processes. The commercial or private laboratory may then compare the gas content measured in some samples with the storage capacity determined from another sample and estimate the degree of saturation of the coal. As explained above, if the measured gas content is less than the storage capacity, the coal is said to be undersaturated with gas, and the laboratory will typically determine the pressure at which the gas content intersects a plot of the storage capacity data. The resulting pressure is typically referred to as the critical desorption pressure (CDP). The CDP is the reservoir pressure at which CBM will start to desorb from the coal with reduction of reservoir pressure, become a gaseous phase, and begin to become capable of production in commercial quantities.\nUnfortunately, the value of CDP determined by the laboratories, too frequently, has been grossly in error from what was ultimately observed when the wells were produced. The present inventor has identified such error in the coring and subsequent laboratory analyses of several of approximately ten wells, analyzed under traditional core analysis using different laboratories. Some analyses have indicated that the reservoirs are saturated at reservoir pressure, yet these reservoirs have not produced any commercial quantities of gas until the reservoir pressure has been drawn down to at least 50 to 60% of the initial reservoir pressure before reaching the CDP. Some of the analyses indicate that the gas contents exceed the storage capacities of the coals at reservoir pressure, something that appears to defy an adequate physical explanation.\nIn summary, coal sampling, coring, and subsequent core analyses as described above may lead to results that are not only time consuming and expensive to obtain, but also they can be highly questionable and frequently inconsistent when used for individualized analysis. For individualized analysis, due to uncertainty, the better use for coal sampling, coring, and core analyses may not come from individual assessments but instead from multiple assessments from which composite isotherms are constructed for a given geological region by averaging of the data and statistically demonstrating the uncertainty. This has been done in the Powder River Basin (PRB) by the United States Bureau of Land Management (BLM) as described in the reference to Crockett and Meyer, 2001. For example, from some 40 samples, the BLM has constructed an averaged synthesized isotherm for samples measured in the PRB representing these 40 samples. Even from such a relatively large number of samples, and ignoring the cost challenges to achieve such data, this effort highlights the challenges in a coal sampling approach because uncertainty in the data still exists. In fact this data shows significantly differing isotherms that represent one standard deviation on either side of the mean curve.\nAnother problem under traditional analysis can, and does, occur in some undersaturated CBM reservoirs when one tries to demonstrate, perhaps through individual testing or small-scale pilots of several adjacent wells, that the well(s) will ultimately produce commercial quantities of CBM. A long and uncertain dewatering period, even under the best of circumstances, may be required before any commercial quantities of CBM are produced. This can lead to long periods of evaluation time. In some areas where there is high permeability and strong aquifer support, such as can be the case in the PRB, one well cannot draw down the pressure sufficiently to ever reach the CDP in any sort of practical or economic time frame. In response to this problem and in an effort to evaluate their leases, most operators have drilled costly (multi-million dollar) multiple-well pilots in an effort to cause interference between wells so that these wells, in combination, can draw the pressure down sufficiently to reach the CDP by exceeding the water influx into the pilot area. Some of these pilots have been successful in the PRB, but some of the pilots have been dewatering for over three years without yet producing commercial quantities of CBM. This dewatering is done at considerable cost of equipment and power to pump wells, at a financial cost of deferred revenues and with the uncertainty that the ultimate resource to be found may not be sufficient to be profitable.\nThe practical challenges of laboratory involvement and sampling difficulty known to exist in a coal sampling-based technique are perhaps highlighted by reference to U.S. Pat. No. 5,785,131 to Gray. Although this reference involves techniques for sensing formation fluids as in gas-oil systems when the fluid itself is of interest, as it relates to the very different aspect of sampling solids containing a substance of interest, it proposes a system for pressurized capture of the samples from entrained particles during drilling. In the reference, these particles of coal or the like are captured and tested on site to avoid some of the mentioned challenges of laboratory testing. As it relates to the solids such as are of interest in the present invention, however, this reference still relies on a capture of the entrained particles and as such it is subject to the uncertainties and other practical limitations discussed above.\nAnother alternative to those techniques based on sampling of the coal itself involves the use of mudlogging during drilling to obtain, at least a qualitative indication of the presence of CBM. Some have even tried to quantify results (Donovan, 2001), but these techniques can leave much to be desired and problems can exist because the system is not usually closed, thus allowing unmeasured gas to escape. Gas-free drilling water is also typically mixed with formation water of different gas content. Further, particle size can need to be estimated, drilling speed recorded, etc. Then, too, results observed by the inventor for the PRB seem to indicate gas contents that are typically far in excess of those observed. Finally, such techniques provide, at best, an estimate for gas content of the coal and do not provide the practical accuracies desired, neither do these techniques provide an estimate for CDP.\nOther than the coal sampling-based techniques mentioned above, efforts (e.g., see Koenig, 1988) have included attempts to determine CDP by producing the well and dropping the pressure, perhaps by bailing or by a pump lowered into the well until gas starts being produced. These techniques can be fraught with problems, some of which are: 1) if a pump is used in the well, its capacity may not be sufficient to draw the well down in a practical testing time frame to determine when gas starts being produced; 2) as the liquid level drops in the well, air may be pulled into the casing from the surface, if the casing is open at the surface, because the pressure in the casing will likely be lower than the atmospheric pressure at the surface, or if the casing is isolated from atmospheric pressure (e.g., shut in) a vacuum may be drawn on the well and a negative gauge pressure (in this document gauge pressure will refer to measurement of pressure above atmospheric pressure where zero gauge pressure would correspond to atmospheric pressure) may result until there is sufficient release of gas from the coal to overcome the vacuum being drawn by the falling liquid level; and 3) by the time the pressure is drawn down sufficiently to see gas production at the surface, the reservoir may already be affected by two-phase flow that may lead to complications in interpretation. This can also produce results inconsistent with later production history."}
{"text": "Field of the Invention\nMethods and apparatuses consistent with the present invention relate to display, image processing control, and more particularly, to a display apparatus, an image processing apparatus and a control method thereof, which can adjust a size of a UI menu item of a user interface according to a user's preference.\nDescription of the Related Art\nGenerally, a display apparatus such as a TV, a mobile phone, a Personal Digital Assistant (PDA), and a Portable Multimedia Player (PMP) displays a user interface (hereinafter referred to as “UI”) to select a UI menu item on a screen. At this time, the user interface may include UI menu items that may be formed in various types. For example, the UI menu items may be arranged in a horizontal and/or vertical line. UI menu items, each of which has a text list, may also be formed. UI menu items, each of which has a label, may be formed in two dimensions. UI menu items, each of which has an elliptic shape or a circular shape to rotate, may be formed.\nIn a conventional art user interface, there are various methods by which a user may select a UI menu item. For example, the user may typically use a moving highlight or a fixed highlight to select a UI menu item by manipulating buttons up and down and/or right and left.\nHowever, depending on the conventional art user interface, a size value of a UI menu item to be displayed on a screen is preset so that the UI menu item having a constant size corresponding to the preset size value is always displayed. Accordingly, if a plurality of UI menu items with constant size are formed, when a current-highlighted UI menu item is far away from a desired UI menu item, highlighting the desired UI menu item may be complicated. In particular, when the UI menu includes a lot of UI menu items, a problem may occur when highlighting the desired UI menu item.\nFurther, according to the conventional art user interface, as described above, the size value of a UI menu item to be displayed on a screen may be preset regardless of the user's preference. Accordingly, the user cannot select a desired UI menu item directly and easily. In particular, when a child or an impaired adult tries to select a desired UI menu item, they may have trouble accessing and controlling their user interface because it takes time to become familiar with how to access and control the user interface without adequate explanation.\nIf an input means such as a touch screen is employed, when the UI menu item includes a preset small size value, it is difficult to exactly select a desired UI menu item without inconvenience. Conversely, even if the UI menu item has a preset size value that is large, the user interface includes limitable UI menu items that can be displayed on a screen at once so that the user interface may not display all desired UI menu items at once."}
{"text": "Containers are constantly packed and unpacked in fulfillment centers that manage movement and storage of items sold on electronic marketplaces. In modern fulfillment centers, automation is heavily used and items are often sorted and moved to different locations by item conveyers. However, containers, such as reusable totes or cardboard shipping boxes, are often loaded with items manually by a human worker. The worker is responsible for packing the container in a way that prevents damage to items during the packing, but also during the transit of the items in the container. Packing containers manually by workers requires extra time for handling of each item and is expensive.\nSome item conveyers transport items directly into containers. In these arrangements, the items typically fall to the bottom of the container and may be damaged by the fall or by an impact from above by a heavy item deposited into the container. For example, when a light bulb is followed by a dictionary along the item conveyer that feeds items into the container, the light bulb would likely get crushed when the dictionary falls on top of the light bulb. One solution is to order the heavier items first along the item conveyer. However, this sortation is often logistically challenging and expensive. For example, additional sortation may require additional conveyers and/or other hardware or sensors."}
{"text": "Myelodysplastic syndromes (MDS) are hematopoietic stem cell malignancies with a rising prevalence owing to the aging of the American population. MDS comprise a group of malignant hematologic disorders associated with impaired erythropoiesis, dysregulated myeloid differentiation and increased risk for acute myeloid leukemia (AML) transformation. The incidence of MDS is increasing with 15,000 to 20,000 new cases each year in the United States and large numbers of patients requiring chronic blood transfusions. Ineffective erythropoiesis remains the principal therapeutic challenge for patients with more indolent subtypes, driven by a complex interplay between genetic abnormalities intrinsic to the MDS clone and senescence dependent inflammatory signals within the bone marrow (BM) microenvironment. Although three agents are approved for the treatment of MDS in the United States (US), lenalidomide (LEN) represents the only targeted therapeutic. Treatment with LEN yields sustained red blood cell transfusion independence accompanied by partial or complete resolution of cytogenetic abnormalities in the majority of patients with a chromosome 5q deletion (del5q), whereas only a minority of patients with non-del5q MDS achieve a meaningful response, infrequently accompanied by cytogenetic improvement. Although responses in patients with del5q MDS are relatively durable, lasting a median of 2.5 years, resistance emerges over time with resumption of transfusion dependence."}
{"text": "It is known to use rotary actuators to control the position of flight control surfaces on an aircraft. An example of a flight control surface is a spoiler, which is positioned on a trailing edge of an aircraft wing. The spoiler can be extended upwardly, from a stowed position flush to the wing to a raised position, to reduce the lift created by the wing. In order to connect the rotary output shaft to the flight control surface, it is common to use an elongate connecting rod.\nIn prior art rotary actuators, it is known to provide some sort of position sensor on the actuator output shaft to determine the angular displacement thereof, which in turn, allows the position of the flight control surface to be determined."}
{"text": "1) Field of the Invention\nThis invention relates to the field of video signal processing, and, in particular, to video signal processing using an architecture having a plurality of parallel execution units.\n2) Background Art\nIt is well known in the prior art to use multiple-instruction multiple-data systems for video signal processing. In a multiple-instruction multiple-data execution of an algorithm each processor of the video signal processor nay be assigned a different block of image data to transform. Because each processor of a multiple-instruction multiple-data system executes its own instruction stream, it is often difficult to determine when individual processors have completed their assigned tasks. Therefore, a software synchronization barrier may be used to prevent any processors from proceeding until all processors in the system reach the same point. However it is sometimes difficult to determine where synchronization barriers are required. If a necessary barrier is omitted by a user then the resulting code may be nondeterministic and re-execution of the code on the same data may yield different results.\nAn alternate architecture known in the prior art is single-instruction multiple-data architecture. Single-instruction, multiple-data is a restricted style of parallel processing lying somewhere between traditional sequential execution and multiple-instruction multiple-data architecture having interconnected collections of independent processors. In the single-instruction, multiple-data model each of the processing elements, or datapaths, of an array of processing elements or datapaths executes the same instruction in lock-step synchronism. Parallelism is obtained by having each datapath perform the same operation on a different set of data. In contrast to the multiple-instruction, multiple-data architecture, only one program must be developed and executed.\nReferring now to FIG. 1, there is shown prior art single-instruction multiple-data architecture 100. A conventional single-instruction multiple-data system, such as architecture 100, comprises a controller 112, a global memory 126 and execution datapaths 118a-n. A respective local memory 120a-n may be provided within each execution datapath 118a-n. Single-instruction multiple-data architecture 100 performs as a family of video signal processors 118a-n united by a single programming model.\nSingle-instruction multiple-data architecture 100 may be scaled to an arbitrary number n of execution datapaths 118a-n provided that all execution datapaths 118a-n synchronously execute the same instructions in parallel. In the optimum case, the throughput of single-instruction multiple-data architecture 100 may theoretically be n times the throughput of a uniprocessor when the n execution datapaths 118a-n operate synchronously with each other. Thus, in the optimum case, the execution time of an application may be reduced in direct proportion to the number n of execution datapaths 118a-n provided within single-instruction multiple-data architecture 100. However, because of overhead in the use of execution datapaths 118a-n, this optimum is never reached.\nArchitecture such as single-instruction multiple-data architecture 100 works best when executing an algorithm which repeats the same sequence of operations on several independent sets of highly parallel data. For example, for a typical image transform in the field of video image processing, there are no data dependencies among the various block transforms. Each block transform may be computed independently of the others.\nThus the same sequence of instructions from instruction memory 124 may be executed in each execution datapath 118a-n. These same instructions are applied to all execution datapaths 118a-n by way of instruction broadcast line 116 and execution may be independent of the data processed in each execution datapath 118a-n. However, this is true only when there are no data-dependent branches in the sequence of instructions. When data-dependent branches occur, the data tested by the branch will, in general, have different values in each datapath. It will therefore be necessary for some datapaths 118a-n to execute the subsequent instruction and other datapaths 118a-n to not execute the subsequent instruction. For example, the program fragment of Table I clips a value v between a lower limit and an upper limit:\nTABLE I ______________________________________ local v; . . . v = expression if (v &gt; UPPER.sub.-- LIMIT) v = UPPER.sub.-- LIMIT; if (v &lt; LOWER.sub.-- LIMIT) v = LOWER.sub.-- LIMIT; ______________________________________\nThe value being clipped, v, is local to each execution datapath 118a-n. Thus, in general, each execution datapath 118a-n of single-instruction multiple-data architecture 100 executing the program fragment of Table I may have a different value for v. In some execution datapaths 118a-n the value of v may exceed the upper limit, and in others v may be below the lower limit. Other execution datapaths 118a-n may have values that are within range. However the execution model of single-instruction multiple-data architecture 100 requires that a single identical instruction sequence be executed in all execution datapaths 118a-n.\nThus some execution datapaths 118a-n may be required to idle while other execution datapaths 118a-n perform the conditional sequence of Table I. Furthermore, even if no execution datapaths 118a-n of single-instruction multiple-data architecture 100 are required to execute the conditional sequence of the program fragment of Table I, all execution datapaths 118a-n would be required to idle during the time of the conditional sequence.\nAnother problem with systems such as prior art single-instruction multiple-data architecture 100 is in the area of input/output processing. This problem arises due to the very large bandwidth requirements associated with this type of architecture. Even in conventional uniprocessor architecture a single block read instruction may take a long period of time to process because memory blocks may comprise a large amount of data in video image processing applications. However, this problem is compounded when there is a block transfer for each enabled execution datapath 118a-n of architecture 100 and datapaths 118a-n must compete for access to global memory 126. For example, arbitration overhead may be very time consuming. This is a problem particularly in video processing because of the large amount of data which must be processed in video applications.\nThe alternative of providing each execution datapath 118a-n with independent access to external memory 126 is impractical for semiconductor implementation. Furthermore, this alternative restricts the programming model so that data is not shared between datapaths 118a-n. Thus further inefficiency results due to the suspension of processing of instructions until all the block reads are completed. This may be seen in the discrete cosine transform image kernel of Table II:\nTABLE II ______________________________________ for (i = 0; i &lt; NUMBEROFBLOCKS; i = i + 4) { k = i + THIS.sub.-- DP.sub.-- NUMBER; read.sub.-- block(original.sub.-- imagek!, temp.sub.-- block); DCT.sub.-- block(temp.sub.-- block); write.sub.-- block(xform.sub.-- imagek!, temp.sub.-- block); }; ______________________________________\nThe read.sub.-- block and write.sub.-- block routines of the instruction sequence of Table II must be suspensive. Each must be completed before the next operation in the kernel is performed. For example, read.sub.-- block fills temp.sub.-- block in local memory 120a-n with all of its local values. These local values are then used by DCT.sub.-- block to perform a discrete cosine transform upon the data in temp.sub.-- block. Execution of the discrete cosine transform must wait for all of the reads of the read.sub.-- block command of all execution datapaths 118a-n to be completed. Only then can the DCT.sub.-- block and write.sub.-- block occur. Thus, by the ordering rules above, read.sub.-- block must be completed before the write.sub.-- block is processed, or the DCT.sub.-- block is executed.\nReferring now to FIG. 2, there is shown processing/memory time line 200. The requirements imposed by the ordering rules within single-instruction multiple data architecture 100 result in the sequentialization of memory transactions and processing as schematically illustrated by processing/memory time line 200. In time line 200, memory read.sub.-- block time segment 202 of execution datapath 118a-n must be completed before processing of DCT.sub.-- block time segment 204 may begin. Processing DCT.sub.-- block time segment 204 must be completed before memory write.sub.-- block time segment 206 may begin. Only when memory write.sub.-- block time segment 206 of a execution datapath 118a-n is complete, can memory read.sub.-- block time segment 208 of a execution datapath 118a-n begin. Execution and access by second execution datapath 118a-n is sequentialized as described for the first.\nThis problem occurs in high performance disk input/output as well. In a typical disk input/output operation an application may require a transfer from disk while continuing to process. When the data from disk are actually needed, the application may synchronize on the completion of the transfer. Often, such an application is designed to be a multibuffered program. In this type of multibuffered program, data from one buffer is processed while the other buffer is being filled or emptied by a concurrent disk transfer. In a well designed system the input/output time is completely hidden. If not, the execution core of single-instruction multiple-data architecture 100 is wait-stated until the data becomes available. This causes further degrading of the performance of the single-instruction multiple-data architecture 100."}
{"text": "A turbocharger uses the high temperature and high pressure energy of the exhaust gas to turn an exhaust turbine at high speed. The exhaust turbine is coaxial with and helps drive the compressor. Along with increasing the rpm of the compressor, it is also capable of raising the pressure in the air intake manifold so that it is greater than atmospheric pressure. From this supercharging pressure, a large volume of intake air can be supplied to the engine.\nFor a vehicle engine with a wide range of rpm speeds, in the medium-to-high speed operating range a sufficient supply of supercharging pressure can be maintained. However in the low speed operating range it is difficult to obtain a sufficient amount of exhaust pressure. Therefore, supercharging pressure cannot be obtained. This leads to an insufficient low speed torque.\nIn this case, a determining factor of the supercharging pressure in the low speed operating range is given by the ratio of the scroll cross-sectional area A, and the radius measured from the center R, that is A/R. If the area A can be made small, even in the low speed operating range where the amount of exhaust gas is small, the rpm can be increased and the rise in supercharging pressure can be quickened.\nIn this way, it has been required to develop the turbocharger so that even in the low speed operating range, a sufficient amount of supercharging pressure can be obtained. However, in this improvement the following problems are uncovered.\nA controlling solenoid valve is necessary to improve the control of supercharging pressure. In the solenoid valve, if there are differences in tolerance from part to part or if changes in tolerance occur with use, the shifts from the target control values caused by this can be corrected by feedback control.\nHowever, since the engine transition time, which is characteristic to the feedback control, is sometimes not in time with the feedback control, the supercharging pressure will shift from the target control value. For example, when the air intake quantity is suddenly increased by stepping on the accelerator pedal, the control value corresponding to the air intake quantity changes and it is necessary to largely make corrections for the shift. Because of this, for a short period of time, it is difficult to follow and respond to the feedback control. During this time the supercharging pressure will differ greatly from the target control value. Therefore, it is required to find a learning amount based on the feedback, correction for use in the succeeding control.\nAlso, when the engine transition time is not in time with the feedback control, it is desirable to use a learning control system that will remove the shift from the target control value and improve the response of the system. The learning control is a system that learns the amount of feedback correction found through the feedback control. Whether or not the proper learning amount is obtained is controlled by the period of the feedback control during which the learning is performed. Improvement in control precision due to learning control cannot be expected if the calculation timing control for the learning is not optimally performed.\nFurther, in the case where a variable geometry means is installed in a turbocharger to control the supercharging pressure, if only the variable geometry means is controlled with the feedback control, the control precision of the controlling range covered by the variable geometry means could be increased. However, if the arrangement of the exhaust bypass valve varies, the control precision of the controlling range covered by the exhaust bypass valve means would decrease.\nIn trying to increase the control precision both of the ranges covered by the exhaust bypass valve means and the variable geometry means simply by performing feedback control with multiple control means, the multiple control methods interfere with each other and proper control cannot be performed. For example, when the variable geometry means shifts from the optimum position to the closing side, the exhaust bypass valve means shifts from the optimum position to the opening side, and as a whole, the supercharging pressure is maintained at the target control value. However, in order to obtain the maximum engine capabilities, the variable geometry means and the exhaust bypass valve means must have optimum positions. If they shift from that position as in the case mentioned above, the turbine capacity becomes small and the engine output drops.\nIn addition, it is effective if overboost control is performed when accelerating suddenly to improve the acceleration capabilities. The overboost control, temporarily raises the target supercharging pressure so that a high supercharging pressure can be obtained improving the engine output. When applying the overboost control to a supercharging pressure control that avoids the interference of multiple control, the control mechanism will switch from the variable geometry means to the exhaust bypass valve means during overboost control. Even if the exhaust bypass valve means is completely closed in the case where the turbine capacity is large, the target supercharging pressure sometimes cannot be obtained during overboost control. That is to say, when suddenly accelerating and overboost control is performed, the target supercharging pressure is raised to the higher target supercharging pressure. In obtaining this higher target supercharging pressure, as previously mentioned, it is necessary that the control duty value in the control through the variable geometry means is large and that the opening of the variable geometry means is made smaller (moves toward the closing side). However, when suddenly accelerating in operating with a large load especially, the operating range to switch the control is attained during overboost control to be switched over to the control through the exhaust bypass valve means. Consequently, the control duty value in the control through the variable geometry means after switching becomes small. This small control duty value is unable to maintain the small opening of the variable geometry means, and adversely the opening of the variable geometry means gets large making it difficult to obtain the higher target supercharging pressure during overboost control."}
{"text": "Conventionally, there are rolling bearing apparatuses to which a rotation detector is installed. There are a passive type and an active type in such rotation detector. The active type can detect a non-rolling state, thereby achieving rotation detection with high precision. The active-type rotation detector includes a pulser ring and a magnetic sensor. The pulser ring is made of a multipolar magnet in which the north pole and the south pole are alternately disposed in a circumferential direction and mounted on the rolling element side out of the inner ring and the outer ring of the rolling bearing. The magnetic sensor is mounted on the non-rolling element side of the rolling bearing opposing the pulser ring. As for the operation, pulse signals are to be outputted from the magnetic sensor according to the rotation of the pulser ring which synchronously rotates with the rolling element and the rotation state of the rolling element is detected by processing the pulse signals. In order to further improve the accuracy of the rotation detection in such active-type rotation detector, it is proposed to reduce each pitch of the magnetic poles of the pulser ring. However, there is a limit in reducing the pitch, which results in limiting the improvement in the accuracy of the rotation detection."}
{"text": "This invention relates to methods of promoting the growth of plants and in particular it relates to methods for increasing the yield of fruit production of plants. In another aspect, the invention relates to increasing the rate of growth of legumes and yields of edible reproduction fruits of said legumes. In still another aspect, the application of effective amounts of the vanadyl compositions applied to plants during growth cycles has been found to promote plant growth and plant fruit yield in peanuts, a legume. In yet another aspect, the invention relates to increase quantity and quality of yields of flowering plants, vegetables and plant growth. Worldwide increasing population increases demands on horticulture efficiency and agricultural crop yields from all sources. The need to improve these efficiencies and yields is ever present especially in view of the loss of productive farm land. Thus, improved plant growth and plant fruit yields and quality are needed along with better utilization of agricultural lands in order to meet the increasing population growth food requirements. Plant growth regulators which can economize the use of fertilizers, nutrients, enzymes, and other growth promoters for various plants are being investigated continuously. Unfortunately, plant growth regulators and specific species of plants along with different generic classes of plants provides the researcher with less than predictable results even when utilizing same or similar compounds successfully applied to other species and generic classes.\nPlant growth regulators can be defined as compounds and preparation which in minute amounts alter the behavior of crop plants or the produce of such plants through physiological (hormonal) tendencies rather than physical action. Plant growth regulators may either accelerate or retard growth, prolong or break a dormant condition, promoting rooting, fruit-set, or increase fruit size or quantity, or effect a growth and or productivity of plants in other ways. Plant growth regulation activity can vary from plant to plant as well as from novel composition to novel composition and specific combinations of the compositions and plants.\nVarious compositions have been studied as micro nutrients and plant growth promoters to increase crop productivity across a wide range of horticulture and agricultural crops. Compositions which include vanadium compounds have been studied and have results both beneficial and detrimental in effect on plants. Vanadium occurs in higher plants at levels usually between 0.2 and 4 ppm. Studies indicate that vanadium may have some effect on plants, including higher order plants which are of agricultural interest. Whether that effect is beneficial or detrimental appears to depend on the form of the vanadium, the type of plant and the method and timing of the application.\nThe importance of many metallic elements as integral constituents of enzymes and electron carriers warrants the attention of many plant physiologists. Aside from their roles as cofactors, activators are regulators in many enzymes' reactions, trace elements are essential to the maintenance of many biochemical processes in plants. Vanadium is one of the trace elements that is widely distributed in plants. Despite the fact that has been shown that vanadium increases growth in lettuce, tomatoes, asparagus, corn, barley and rice and the like, there is still no conclusive evidence that it is an essential element for higher plants. It does not meet the criteria of essentiality to the plant. However, vanadium on the other hand appears to be essential for many microorganisms. Continuing studies are being undertaken presently as they have in the past to investigate the distribution of vanadium and its effects on growth of agriculturally useful plants. Vanadium has stimulated growth and mazed plants at levels of 0.25 ppm in nutrient solutions. However, vanadium had no effect on lettuce or tomato plants at levels of 0.05 ppm. Lauchili et al. 15B ENCYCLOPEDIA OF PLANT PHYSIOLOGY 723-26 (1983). A study of the effect of a foliar spray of a vanadyl sulfate solution on leaf growth of sugar beet plants indicated that it decreased leaf growth, but that it increased the amount of reducing sugar in the roots of the sugar beet. Singh et al. 44 PLANT PHYSIOLOGY 1321-27 (1969). The use of vanadium (in the form of vanadyl lactate at concentrations of 10.sup.-3 to 10.sup.-6 molar) as a fertilizer has produced an increase in foliage yield of some higher plants. Kerr et al. Monograph 11 BRITISH PLANT GROWTH REGULATOR GROUP 103-21 (1984).\nMore recently vanadyl composition have been used in promoting the lint yield of fiber producing plants such as cotton by heating the plant during growth periods of the plant with an amount of a vanadyl salt of a carboxycilic acid. These compositions are found in U.S. Pat. No. 5, 186,738 hereby incorporated by reference; however the '738 patent reports that the vanadyl composition proved ineffective for promoting the growth of soybeans, a legume. Again illustrating the variations of experience with vanadyl compositiions in achieving regulation activity ie. plant growth and plant fruit growth enhancement from plant to plant. Thus, a continuing need exists for development of effective vanadium composition methods, timing and conditions of application in order to achieve enchanced edible fruit production."}
{"text": "The ultimate goal of a robotics system is of course to perform some useful function in place of its human counterpart. Benefits typically associated with the installation of fixed-location industrial robots are improved effectiveness, higher quality, reductions in manpower, greater efficiency, reliability and cost savings. Additional benefits include the ability to perform tasks for which humans are incapable and the removal of humans from dangerous or life-threatening scenarios. The concept of mobility has always suggested an additional range of applications beyond that of the typical factory floor, wherein free-roaming robots moved about with an added versatility that brought even greater returns. In practice, however, the realization of this dream has been slow in coming.\nOne of the most significant technological hurdles impeding the widespread introduction of mobile robotics systems arises from the need for a mobile platform to successfully interact with the physical objects and entities in its environment. The robot must be able to navigate from a known position to a desired new location and orientation, at the same time avoiding any contact with fixed or moving objects while enroute. A robust navigational scheme that preserves the validity of a world model for free-roaming platforms has remained an elusive research goal, and for this reason many proposed applications of autonomous mobile robots are yet to be implemented.\nThe simplest form of autonomous control is sometimes termed guidepath control and involves a navigational control loop which reacts (in a reflexive manner) to the sensed position of some external guiding reference. The intent is to free a human operator from the requirement of steering the moving platform. For example, encoded reflective stripes might be applied to the floor of the robot's environment. The robot would then be equipped with stripe detecting/decoding for determining the robot's position in its environment as provided on the encoded stripes. Such automated guided vehicles (AGVs) have found extensive use in factories and warehouses for material transfer, in modern office scenarios for material and mail pickup and delivery, and in hospitals for delivery of meals and supplies to nursing stations, to name but a few.\nAdvantages of guidepath control are seen primarily in the improved efficiency and reduction of manpower since an operator is no longer required to guide the vehicle. Large numbers of AGVs can operate simultaneously in a plant or warehouse without getting lost or disoriented. The AGVs are typically scheduled and controlled by a central computer which monitors overall system operation and vehicle flow. Communication with individual vehicles can be via RF links or directional near-infrared modulated light beams, or other means. However, the fundamental disadvantage of guidepath control is the lack of flexibility in the system. A vehicle cannot be commanded to go to a new location unless the guidepath is first modified. This is a significant drawback in the event of changes to product flow lines in assembly plants, or in the case of a security robot which must investigate a potential break-in at a designated remote location.\nThus, truly autonomous control implies the ability of a mobile platform to travel anywhere so desired, subject to nominal considerations of terrain. Many potential applications await an indoor robot that could move in a purposeful fashion from room to room without following a set guidepath, with the intelligence to avoid objects and, if necessary, choose alternative routes of its own planning. To do this, specialized sensors must be coupled with some type of \"world modeling\" capability that represents the relative/absolute locations of objects detected by these sensors. In this way, a mobile platform can be provided with sufficient awareness of its surroundings to allow it to move about in a realistic fashion, i.e., a path not forever dictated by a guidepath stripe.\nThe accuracy of this model, which is constructed and refined in a continuous fashion as the robot moves about its workspace, is directly dependent throughout this process upon the validity of the robot's perceived location and orientation. Accumulated dead-reckoning errors can quickly render the information entered into the model invalid in that the associated geographical reference point for data acquired relative to the robot's position is incorrect. As the accuracy of the model degrades, the ability of the robot to successfully navigate and avoid collisions diminishes rapidly, until it fails altogether. One solution to this problem is to provide navigation landmarks within the robot's environment for use by the robot in periodic (absolute) positional updates. The concept of using existing interior doorways as navigation landmarks for a mobile robotics system has always been appealing, in that no modifications to the surrounding environment are required. The robot by necessity must travel through a doorway to enter an adjoining space. If in so doing the system could obtain an accurate positional update, then such would indeed represent an elegant solution to the problem of cumulative dead-reckoning errors.\nThus, the need exists for a navigational referencing system for a mobile robot that can derive its own updated positional information as it moves through its environment. Accordingly, an object of the present invention is to provide a navigational referencing method and system for a mobile robot that derives x-y position and angular orientation of the robot within a world model of the robot's environment as the robot traverses doorway openings. Another object of the present invention is to provide a navigational referencing method and system for a mobile robot that derives a relative x-y position and angular orientation of the robot without any modifications to the robot's environment. Yet another object of the present invention is to provide a method and system of navigational referencing for a mobile robot that provides for sufficient updates to the x-y position and angular orientation of the robot within a world model of the robot's environment thereby avoiding the accumulated effect of navigational dead-reckoning errors."}
{"text": "Processes executed in a computer system often have tasks with different processing demands and priorities. In order to operate as desired, a process may expend significant overhead to ensure that suitable resources are allocated to tasks based on the demands and priorities of the tasks. This overhead may include the use of one or more schedulers that schedule tasks of the process for execution in the computer system. The execution of tasks with schedulers may vary based on the processing capabilities and/or processing demands of the computer system at any given time. As a result, resource allocations to schedulers may not be optimal for some processing conditions."}
{"text": "1. Field of the Invention\nThe present invention generally relates to noise reduction methods and image decoding devices, and particularly relates to a method of reducing quantization noise generated during a decoding process of coded image data and to a device for decoding the coded image data based on the method.\n2. Description of the Related Art\nWhen transmitting, recording, or reproducing various signals such as video signals, audio signals, etc., as digital signals, techniques for compressing and decompressing information are generally employed to reduce the amount of information, i.e., the number of bits. For example, if a linear quantization (uniform quantization) which represents each sample value with a value selected from evenly divided signal levels is used for digitizing video signals;, audio signals, and the like without any compression technique, the amount of information to be transmitted or recorded/reproduced becomes prohibitively large.\nThus, in the fields of broadcasting, communication, and information recording/reproducing, characteristics of human visual perception and auditory perception have been utilized in such compression techniques. For example, the human perception is sensitive to changes in signal levels when a signal has a little variation, while it is not so sensitive to the changes when a signal has a strong fluctuation. Such characteristics can be utilized to reduce the amount of information for each sample value. Also, a number of technologies for compressing information have been employed to bring about an advance in the practical use of highly efficient compression techniques.\nFor example, the amount of information contained in a one-hour moving picture having an image quality similar to those reproduced by VHS-type VTR devices is about 109 Gbits. Also, 360 Gbits more or less are contained in such a one-hour moving picture with an image quality comparable to reception images of NTSC color-television sets. Thus, an effort to develop highly efficient compression techniques is also directed to application studies aimed at transmitting or recording/reproducing such a large amount of information by means of current transmission lines or recording media.\nHighly efficient compression methods which have been proposed as practical methods for image-information application typically combine three different compression techniques to reduce the amount of information. The first technique reduces the amount of information by utilizing correlation within an image frame (compression utilizing spatial correlation), which takes advantage of the fact that there is high correlation between adjacent pixels in natural images. The second technique reduces the amount of information by utilizing correlation between image frames arranged in time (compression utilizing temporal correlation). The third technique reduces the amount of information by utilizing a different probability of appearance of each code.\nAs techniques for compressing image information by utilizing correlation within an image frame (the first technique), a variety of techniques have been proposed. In recent years, orthogonal transformation such as the K-L (Karhunen-Loeve) transform, the Discrete Cosine Transform (DCT), the discrete Fourier transform, and the Walsh-Hadamard transform have often been employed.\nFor example, highly efficient coding methods for image information proposed by the MPEG (moving picture coding expert group) which has been established under the ISO (international standardization organization) employ two-dimensional DCT. These highly efficient coding methods (MPEG1 and MPEG2) combine intra-frame coding and inter-frame coding to realize highly efficient coding of moving-picture information while employing motion compensation prediction and inter-frame prediction. The orthogonal transformation is generally applied to blocks which are generated by dividing an image into unit blocks having a predetermined block size (M.times.N). In MPEG1 and MPEG2, a block having an 8-pixel-by-8-pixel block size is defined as a unit block.\nM.times.N orthogonal transform coefficients which are obtained by orthogonally transforming the unit block (e.g., 64 DCT transform coefficients in MPEG1 and MPEG2) are then quantized by using block-quantization step sizes (intervals for quantization). The block-quantizatoin step sizes are defined for each predetermined-size area including at least one unit block. In MPEG1 and MPEG2, for example, this predetermined-size area is called a macro block, which consists of a block of 16.times.16 pixels for a luminance signal Y and a block of 8.times.8 pixels for each of color signals Cr and Cb. In detail, the block-quantization step sizes are represented as [{a quantization characteristic of a macro block (a quantization scale of a macro block) QS}.times.quantization matrix (8.times.8)]. Here, the quantization characteristic of a macro block changes from macro block to macro block.\nThe orthogonal transform coefficients (e.g., DCT coefficients) which are quantized based on the block-quantization step sizes are separated into a direct current component (DC component) and alternating current components (AC components). The direct current component of the orthogonal transform coefficients is subjected to differential coding, and the alternating current components of the orthogonal transform coefficients are subjected to entropy coding after a zigzag scan. Here, the entropy coding is an information compression technique using a variable-length coding scheme which utilizes a different probability of appearance of each code such as in the Huffman coding.\nTransformed and coded image data is transmitted as a bit stream (a series of bits). A decoding operation on the transformed and coded image data is carried out in a reversed manner to the coding operation described above so as to generate an output image. However, when the quantization process is included in the entire coding process, unavoidable quantization errors result in quantization noise appearing in the output image. Thus, when the complexity of an image subjected to the coding process contains a larger amount of information than capacity of a transmission rate, the quantization noise will substantially degrade image quality.\nIn general, the quantization errors in low-frequency components result in block distortions in the output image, by which there appears to be no correlation between each block of the output image. Also, the quantization errors in high-frequency components generate mosquito noise around edges, which is a distortion having a ringing appearance in the output image.\nThe quantization errors appearing in the output image are especially conspicuous where image levels are generally flat. When a small amount of the quantization noise is added at a point where a change in a video signal level has frequency components from low frequencies to high frequencies, the noise is difficult to be visually detected because of characteristics of visual perception. However, when a small amount of noise having high frequency components is added at a point where a change in the video signal has only low frequency components, the noise is easy to detect. Of course, when a large amount of noise is added, the noise is detected as coding degradation irrespective of the frequency components of the noise.\nIn order to obviate the problem of an image-quality degradation caused by the quantization noise, Japanese Laid-Open Patent Application No. 4-2275 discloses a technique for reducing the quantization errors in low frequency components which appear as block distortion in an image. The technique detects an activity value (the amount of high frequency components contained in a given block of image data) for a given block. Based on the activity value, pixels near the borders of the given block are processed by a low-pass filter, and, then, random noise is added to these pixels processed by the low-pass filter.\nHowever, this technique has a problem when the image degradation is substantial as in the case where an image to be subjected to highly efficient coding is relatively complex in comparison with a transmission rate. In such a case, the quantization noise causing the image degradation is mistakenly detected as part of the activity value of the decoded image. Thus, even if a block is of a low activity value, this block might be mistakenly judged as having a high activity value. That is, it is difficult to determine whether the activity value detected in the decoded image reflects a real activity value inherent in the image data or reflects the quantization noise sneaking into the image data. As a result, the technique described above cannot effectively reduce the block distortion generated during the decoding process of the coded image data.\nAccordingly, there is a need for a method which can effectively reduce the quantization noise generated during a decoding process of coded image data, and for a device for decoding the coded image data based on this method."}
{"text": "The present invention relates to article processing systems and, more particularly, to a conveyor system for moving an article along a conveying path to successive processing stations.\nArticle processing systems commonly employ an elevator/conveyor system that moves articles to be treated or processed to work stations positioned along a conveyor path. An elevator mechanism is utilized to raise and lower the articles to and from selected work stations, e.g., processing tanks, along the path. In this manner, the articles are successively moved from a loading/unloading station to and through a succession of work stations and returned to the loading/unloading station. These conventional processing systems are well suited for processing a plurality of articles according to a predetermined procedure. However, such conventional processing systems are often expensive to produce and costly to operate and maintain because of the relative complexity of their design.\nAnother type of system employs a continuous conveyor with fixed successive declined and inclined sections for immersing an object in a processing tank and subsequently removing the object from the tank. The degree of decline and incline necessitates relatively large volume processing tanks and extended conveying paths. Further, the time between processing steps tends to be excessive in some cases.\nAccordingly, there is a need for a conveyor system that enables efficient mass production processing and minimizes the costs associated with design, manufacture, operation, and maintenance of the conveyor system."}
{"text": "The physical uplink channels of the long term evolution (LTE for short) system include a physical random access channel (PRACH for short), a physical uplink shared channel (PUSCH for short) and a physical uplink control channel (PUCCH for short). Wherein, the PUSCH has two different cyclic prefix (CP for short) lengths which are respectively a normal cyclic prefix (Normal CP for short) and an extended cyclic prefix (Extended CP for short). Each sending subframe of the PUSCH is composed of two time slots. For different cyclic prefix lengths, the location of the demodulation reference signal (DMRS for short) in the subframe will be different. FIG. 1 is a schematic diagram of a time domain location of a demodulation reference signal according to the related art. As shown in FIG. 1, each subframe contains two DMRS symbols. FIG. 1 (a) is a schematic diagram of the time domain location of the DMRS when adopting the normal cyclic prefix, each subframe contains 14 orthogonal frequency division multiplexing (OFDM for short) symbols, and the 14 OFDM symbols include the DMRS symbols, wherein, the OFDM symbols represent the time domain location of one subframe; FIG. 1 (b) is a schematic diagram of the time domain location of the DMRS when adopting the extended cyclic prefix, and each subframe contains the OFDM symbols of 12 time domains.\nIn the LTE system, a physical downlink control channel (PDCCH for short) is used to bear the uplink and downlink scheduled information, and the uplink power control information. A base station (e-Node-B, eNB for short) can configure the user equipment (UE for short) through the downlink control information, or the user equipment accepts the configuration from the higher layers, which is also called as configuring the UE through the high layers signaling. The format of the downlink control information (DCI for short) includes DCI format 0, 1, 1A, 1B, 1C, 1D, 2, 2A, 3 and 3A, etc., wherein,\nthe DCI format 0 is used to indicate scheduling of the PUSCH;\nthe DCI format 1, 1A, 1B, 1C and 1D are used for different transmission modes of a physical downlink shared channel (PDSCH for short) of a single transport block;\nthe DCI format 2 and 2A are used for different transmission modes of space division multiplexing of the downlink PDSCH;\nthe DCI format 3 and 3A are used for transmission of power control instructions of the PUCCH and the PUSCH.\nThe transport block size of the above-mentioned DCI format 0, 1A, 3 and 3A are same, wherein, the DCI format 0 and 1A adopts 1 bit to distinguish the format.\nThe format of the DCI format 3 is as follows:                transmission power control command 1, transmission power control command 2, . . . , transmission power control command N,        \nwherein,\n      N    =          ⌊                        L                      format            ⁢                                                  ⁢            0                          2            ⌋        ,Lformat 0=format 0 plus a bit number before the cyclical redundancy check (CRC for short) (including additional padding bit(s)), and the parameter tpc-Index given by the high layers determines the transmission power control command (TPC command) of the given UE.\nIf\n            ⌊                        L                      format            ⁢                                                  ⁢            0                          2            ⌋        <                  L                  format          ⁢                                          ⁢          0                    2        ,the DCI format 3 will be added 1 bit ‘0’.\nThe process of a blind detection of the PDCCH in the LTE system is described as follows briefly, the control channel element (CCE for short) is a minimum element bearing the PDCCH resource, and the control area is composed of a series of CCEs.\nThe blind detection range of the PDCCH is defined by a search space, and the search space is divided into a public search space and an UE dedicated search space. The search space Sk(L) is defined as:L·{(Yk+m)mod └NCCE,k/L┘}+i, \nwherein, L is the aggregation grade of the CCE, and Lε{1, 2, 4, 8}; for the public search space, Yk=0, i.e., searching from CCE=0˜15; and for the UE dedicated search space, Yk=(A·Yk-1)mod D, Y−1=nRNTI≠0, A=39827, D=65537, k=└ns/2┘, ns represents the time slot number 0˜19. i=0,L, L−1, m=0,L, M(L)−1, M(L) is the number of PDCCH candidates after L is given in the search space.\nWherein, the nRNTI represents the radio network temporary identifier (RNTI for short), and nRNTI corresponds to one of the following radio network temporary identifiers:\nsystem information-RNTI (SI-RNTI for short),\nrandom access-RNTI (RA-RNTI for short),\npaging-RNTI (P-RNTI for short),\ncell-RNTI (C-RNTI for short),\nsemi-persistent scheduling RNTI (SPS-RNTI for short), and\ntemporary cell-RNTI (Temporary C-RNTI).\nWhich kind of RNTI the nRNTI selects specifically is configured by the high layers signaling, and the specific value is also specified by the corresponding signaling and data. The value of the RNTI refers to the following Table 1. The search space defined according to the aggregation grade is shown in Table 2. When the UE is blind detected, the detection is performed according to the DCI format corresponding to the transmission mode of the downlink. The 16-bit CRC of each downlink control information DCI is scrambled by using the above-mentioned RNTI. Different UEs can configure different RNTIs to perform the scrambling to the CRC, thus it can distinguish the DCI of different UEs.\nTABLE 1RNTI valueValue (Hex)Frequency DivisionTime DivisionDuplexing (FDD)Duplexing (TDD)RNTI0000-00090000-003BRadom access RNTI (RA-RNTI)000A-FFF2003C-FFF2C-RNTI, Semi-PersistentScheduling C-RNTI, TemporaryC-RNTI, TPC-PUCCH-RNTIand TPC-PUSCH-RNTIFFF3-FFFCReservedFFFEP-RNTIFFFFSI-RNTI\nTABLE 2the PDCCH candidates monitored by UESearch space Sk(L)The number ofAggregationPDCCH candidatesTypegrade LSize [CCEs]M(L)UE166dedicated21264828162Public41648162\nIn the LTE system, the multiplexing process of the DCI is shown in FIG. 2. Every DCI corresponds to one medium access control identifier (MAC id), i.e., corresponds to one RNTI. The original information bit of the DCI is added with the CRC scrambled by the RNTI, and then performed channel coding and rate matching, thus a plurality of DCIs of the PDCCH are multiplexed together. System information (SI for short) performs the resource allocation through the DCI format 1A/1C. The blind detection of the SI is only performed in the public search space, and the CRC of the DCI of the SI is scrambled by adopting the unique SI-RNTI.\nThe broadcast information of the LTE system is divided into a master information block (MIB for short) and a system information block (SIB for short), wherein, the MIB is transmitted on the physical broadcast channel (PBCH for short), and the SIB is transmitted on the PDSCH (also called scheduled information (SI)).\nThe SRS is a signal used for measuring radio channel state information (CSI for short) between the user equipment and the base station. In the long term evolution system, the UE sends an uplink SRS in the last data symbol of the sent subframe regularly according to parameters, such as the bandwidth, the frequency domain location, the sequence cyclic shift, the period and the subframe offset, etc., instructed by the eNB. The eNB judges the UE uplink CSI according to the received SRS, and performs operations, such as frequency domain selection scheduling and closed-loop power control, etc., according to the obtained CSI.\nIn the LTE system, an SRS sequence sent by the UE is obtained by performing the cyclic shift α in time domain to one root sequence ru,v(n). Different SRS sequences can be obtained by performing different cyclic shifts α to the same root sequence, and these obtained SRS sequences are mutual orthogonal. Therefore, these SRS sequences can be allocated to different UEs for using, to realize the code division multiple access among the UEs. In the LTE system, the SRS sequence defines 8 cyclic shifts α, which is provided through the following formula (1):\n                    α        =                  2          ⁢          π          ⁢                                          ⁢                                    n              SRS              cs                        8                                              Formula        ⁢                                  ⁢                  (          1          )                    \nwherein, nSRScs is indicated by signaling of 3-bit, as 0, 1, 2, 3, 4, 5, 6 and 7 respectively. That is to say, with the same time-frequency resource, the UE in the cell has 8 available code resources, and the eNB can configure at most 8 UEs to send the SRS at the same time on the same time-frequency resource. Formula (1) can be regarded as dividing the SRS sequence into 8 pieces with same interval in the time domain, however, since the length of the SRS sequence is a multiple of 12, so the minimum length of the SRS sequence is 24.\nIn the LTE system, the frequency domain bandwidth of the SRS adopts a tree structure to configure. Every SRS bandwidth configuration corresponds to one tree structure, and the SRS-Bandwidth on the highest layer (or called the first layer) corresponds to the maximum SRS bandwidth of the SRS bandwidth configuration, or is called as the SRS bandwidth range. The UE calculates and obtains its own SRS bandwidth according to the signaling indication of the base station, and then determines a frequency domain initial position for sending the SRS by itself according to an upper layer signaling frequency domain location nRRC sent by the eNB. FIG. 3 is a schematic diagram of a frequency domain initial position of a UE allocated with different nRRC sending the SRS in the related art. As shown in FIG. 3, the UE allocated with different nRRC will send the SRS in different areas of the SRS bandwidth of the cell, wherein, the UE1 determines the frequency initial position to send the SRS according to nRRC=0, the UE2 determines the frequency initial position to send the SRS according to nRRC=3, the UE3 determines the frequency initial position to send the SRS according to nRRC=4, and the UE4 determines the frequency initial position to send the SRS according to nRRC=6.\nThe sequence used by the SRS is selected from a demodulation pilot frequency sequence group. When the SRS bandwidth of the UE is 4 resource blocks (RB for short), it uses a computer generated (CG for short) sequence with the length being 2 RBs; and when the SRS bandwidth of the UE is larger than 4 RBs, it uses a Zadoff-Chu sequence with the corresponding length.\nIn addition, in the same SRS bandwidth, the sub-carrier of the SRS is placed with intervals, that is to say, sending of the SRS adopts a comb structure. The number of the frequency comb in the LTE system is 2, which also corresponds to the repetition factor (RPF for short) of the time domain of 2. FIG. 4 is a schematic diagram of a comb structure of the SRS in the related art. As shown in FIG. 4, when each UE sends the SRS, only one of the two frequency combs is used, comb=0 or comb=1. In this way, the UE, according to the location indication of the frequency comb (comb=0 or comb=1) of 1-bit upper layer signaling, only uses the sub-carrier with frequency domain index being the even number or odd number to send the SRS. This comb structure allows more UEs to send the SRS in the same SRS bandwidth.\nIn the same SRS bandwidth, a plurality of UEs can use different cyclic shifts in the same frequency comb, and then sends the SRS through the code division multiplexing; and also two UEs send the SRS in different frequency combs through the frequency division multiplexing. For example, in the LTE system, for the UE sending the SRS in a certain SRS bandwidth (4 RBs), cyclic shifts able to be used by the UE are 8 and frequency combs able to be used by the UE are 2, thus the UE has 16 resources that can be used for sending the SRS, that is to say, in that SRS bandwidth, at most 16 SRSs can be sent at the same time. Since the LTE system does not support single user multiple input multiple output (SU-MIMO for short), the UE can only have one antenna to send the SRS at every moment; so one UE only needs one SRS resource. Therefore, in the above-mentioned SRS bandwidth, the system at most can multiplex 16 UEs at the same time.\nThe LTE-Advanced (LTE-A for short) system is the next generation evolved system of the LTE system, supports the SU-MIMO in the uplink, and can use at most 4 antennas as uplink transmitting antennas. That is to say, the UE can simultaneously send the SRS on a plurality of antennas at the same moment, while the eNB needs to estimate the state on each channel according to the SRS received on each antenna.\nIn the case of the LTE-A carrier aggregation, a variety of carrier types are introduced. The LTE-A carrier types can be divided into three types: backwards compatible carrier, non-backwards compatible carrier and extension carrier.\nThe extension carrier has two kinds of meanings: 1) as a part of component carrier (CC for short); 2) as an independent component carrier. The extension carrier cannot work alone, and must be a part of a group of component carrier set; and at least one of the component carriers in the set can work alone. The extension carrier is invisible for the LTE UEs.\nIn order to design simply, and considering various possible application scenarios, the extension carrier is most likely configured to have no PDCCH. Then the DCI corresponding to the system information of the extension carrier needs to be sent on other component carriers. Besides, the LTE-A has also introduced the concept of resident carrier, that is, the carrier accessed by the UE initially, and after the access succeeds, it can reconfigure the resident carrier for the UE through the high layers signaling, to guarantee load balancing.\nIn the case of the LTE-A carrier aggregation, a PDCCH component carrier set (PDCCH CC set) is defined, and the UE needs to perform the blind detection in the PDCCH CC set; a downlink component carrier set (DL CC set for short) is defined also, and the PDSCH of the UE can be sent on any CC in the DL CC set. In the case of the LTE-A carrier aggregation, cross carrier scheduling is allowed, that is, the PDCCH on a certain component carrier can schedule the PDSCH or the PUSCH on a plurality of component carriers.\nIt is proposed in the existing research of the LTE-A that: in the uplink communication, the non-precoding (i.e., antenna-dedicated) SRS should be used, while the precoding is performed to the DMRS of the PUSCH. The base station, by receiving the non-precoding SRS, can estimate the uplink original CSI, but the base station cannot estimate the uplink original CSI through the precoding DMRS. At the moment, when the UE uses multi-antenna to send the non-precoding SRS, the SRS resource required by each UE will increase, which causes the number of the UEs which can be multiplexed in the system at the same time to decrease. In addition, except keeping the original periodic of the LTE to send SRS, in order to improve the utilization rate of the SRS resource and enhance the flexibility of resource scheduling, the UE aperiodically sending the SRS can also be configured through the downlink control information or the high layers signaling. Therefore, there are the periodic SRS and the aperiodic SRS in the LTE-A system, and how to reasonably design the downlink control information or the high layers signaling to configure the SRS resource, to realize aperiodically sending the SRS effectively and timely, save the signaling overhead and reduce the complexity of the UE blind detection at the same time, is a problem to be solved."}
{"text": "This invention relates to telemanipulation using telepresence, and more particularly to applications of telemanipulation to laparoscopic surgery.\nA telemanipulation system allows an operator to manipulate objects located in a workspace from a remote control operator's station. For example, in a laparoscopic abdominal surgery procedure, the patient's abdomen is insufflated with gas, and cannulas are passed through small incisions to provide entry ports for laparoscopic surgical instruments. Laparoscopic surgical instruments include an image capture means for viewing the surgical field and working tools, such as forceps or scissors. The working tools are similar to those used in open surgery, except that the working end of each tool is separated from its handle by an extension tube. The surgeon performs surgery by sliding the instruments through the cannulas and manipulating them inside the abdomen while referencing a displayed image of the interior of the abdomen. Surgery by telepresence, that is, from a remote location by means of remote control of the surgical instruments, is a next step. A surgeon is ideally able to perform surgery through telepresence, which, unlike other techniques of remote manipulation, gives the surgeon the feeling that he is in direct control of the instruments, even though he only has remote control of the instruments and view via the displayed image.\nThe effectiveness of telepresence derives in great measure from the illusion that the remote manipulators are perceived by the operator of the system to be emerging from the hand control devices located at the remote operator's station. If the image capture means, such as a camera or laparoscope, are placed in a position with respect to the manipulators that differs significantly from the anthropomorphic relationship of the eyes and hands, the manipulators will appear to be located away from the operator's hand controls. This will cause the manipulators to move in an awkward manner relative to the viewing position, inhibiting the operator's ability to control them with dexterity and rapidity. However, it is often unavoidable in applications such as laparoscopic surgery to move the laparoscope in order to obtain the best possible image of the abdominal cavity.\nThus, a technique is needed for providing to the operator the sense of direct hand control of the remote manipulator, even in the presence of a substantially displaced imaging device, such that the operator feels as if he is viewing the workspace in true presence."}
{"text": "The present application relates to systems, devices and methods for memory access operations involving phase change memory units.\nNote that the points discussed below may reflect the hindsight gained from the disclosed inventions, and are not necessarily admitted to be prior art.\nPhase change memory (“PCM”) is a relatively new nonvolatile memory technology, which is very different from any other kind of nonvolatile memory. First, the fundamental principles of operation, at the smallest scale, are different: no other kind of solid-state memory uses a reversible PHYSICAL change to store data. Second, in order to achieve that permanent physical change, an array of PCM cells has to allow read, set, and reset operations which are all very different from each other. The electrical requirements of the read, set, and reset operations make the peripheral circuit operations of a PCM very different from those of other nonvolatile memories. Obviously some functions, such address decoding and bus interface, can be the same, but the closest-in parts of the periphery, which perform set, reset, and read operations on an array or subarray, must satisfy some unique requirements.\nThe physical state of a PCM cell's memory material is detected as resistance. For each selected cell, its bitline is set to a known voltage, and the cell's access transistor is turned on (by the appropriate wordline). If the cell is in its low-resistance state, it will sink a significant current from the bit line; otherwise, it will not.\nSet and Reset operations are more complicated. Both involve heat. As discussed below, a “set” operation induces the memory material to recrystallize into its low-resistance (polycrystalline) state, while a “reset” operation anneals the memory material into its high-resistance (amorphous) state.\nWrite operations (Set and Reset) normally have more time budget than read operations. In read mode a commercial PCM memory should be competitive with the access speed (and latency if possible) of a standard DRAM. If this degree of read speed can be achieved, PCM becomes very attractive for many applications.\nThe phase change material is typically a chalcogenide glass, using amorphous and crystalline (or polycrystalline) phase states to represent bit states.\nA complete PCM cell can include, for example: a top electrode (connected to the bit line), a phase change material (e.g. a chalcogenide glass), a conductive pillar which reaches down from the bottom of the phase change material, an access transistor (gated by a word line), and a bottom connection to ground. The phase change material can extend over multiple cells (or over the whole array), but the access transistors are laterally isolated from each other by a dielectric.\nFIG. 2A shows an example of a PCM element 2010. A top electrode 2020 overlies a phase change material 2030, e.g. a chalcogenide glass. Note that material 2030 also includes a mushroom-shaped annealed zone (portion) 2070 within it. (The annealed zone 2070 may or may not be present, depending on what data has been stored in this particular location.) The annealed zone 2070, if present, has a much higher resistivity than the other (crystalline or polycrystalline) parts of the material 2030.\nA conductive pillar 2050 connects the material 2030 to a bottom electrode 2040. In this example, no selection device is shown, though in practice, an access transistor would normally be connected in series with the phase change material. The pillar 2050 is embedded in an insulator layer 2060.\nWhen voltage is applied between the top 2020 and bottom 2040 electrodes, the voltage drop will appear across the high-resistivity zone 2070 (if present). If sufficient voltage is applied, breakdown will occur across the high-resistivity zone. In this state the material will become very conductive, with large populations of mobile carriers. The material will therefore pass current, and current crowding can occur near the top of the pillar 2050. The voltage which initiates this conduction is referred to as the “snapback” voltage, and FIG. 2C shows why.\nFIG. 2C shows an example of instantaneous I-V curves for a device like that of FIG. 2A, in two different states. Three zones of operation are marked.\nIn the zone 2200 marked “READ,” the device will act either as a resistor or as an open (perhaps with some leakage). A small applied voltage will result in a state-dependent difference in current, which can be detected.\nHowever, the curve with open circles, corresponding to the amorphous state of the device, shows some more complex behaviors. The two curves show behaviors under conditions of higher voltage and higher current.\nIf the voltage reaches the threshold voltage Vth, current increases dramatically without any increase in voltage. (This occurs when breakdown occurs, so the phase-change material suddenly has a large population of mobile carriers.) Further increases in applied voltage above Vth result in further increases in current; note that this upper branch of the curve with hollow circles shows a lower resistance than the curve with solid squares.\nIf the applied voltage is stepped up to reach the zone 2150, the behavior of the cell is now independent of its previous state.\nWhen relatively large currents are applied, localized heating will occur at the top of the pillar 2050, due to the relatively high current density. Current densities with typical dimensions can be in the range of tens of millions of Amperes per square cm. This is enough to produce significant localized heating within the phase-change material.\nThis localized heating is used to change the state of the phase-change material, as shown in FIG. 2B. If maximum current is applied in a very brief pulse 2100 and then abruptly stopped, the material will tend to quench into an amorphous high-resistivity condition; if the phase-change material is cooled more gradually and/or not heated as high as zone 2150, the material can recrystallize into a low-resistivity condition. Conversion to the high-resistance state is normally referred to as “Reset”, and conversion to the low-resistance state is normally referred to as “Set” (operation 2080). Note that, in this example, the Set pulse has a tail where current is reduced fairly gradually, but the Reset pulse does not. The duration of the Set pulse is also much longer than that of the Reset pulse, e.g. tens of microseconds versus hundreds of nanoseconds.\nFIG. 2D shows an example of temperature versus resistivity for various PCM materials. It can be seen that each curve has a notable resistivity drop 2210 at some particular temperature. These resistivity drops correspond to phase change to a crystalline (or polysilicon) state. If the material is cooled gradually, it remains in the low resistivity state after cooling.\nIn a single-bit PCM, as described above, only two phases are distinguished: either the cell does or does not have a significant high-resistivity “mushroom cap” 2070. However, it is also possible to distinguish between different states of the mushroom cap 2070, and thereby store more than one bit per cell.\nFIG. 2E shows an equivalent circuit for an “upside down” PCM cell 2010. In this example the pass transistor 2240 is gated by Wordline 2230, and is connected between the phase-change material 2250 and the bitline 2220. (Instead, it is somewhat preferable to connect this transistor between ground and the phase-change material.\nFIG. 2F shows another example of a PCM cell 2010. A bitline 2220 is connected to the top electrode 2020 of the phase-change material 2250, and transistor 2240 which is connected to the bottom electrode 2030 of the PCM element. (The wordline 2230 which gates the vertical transistor 2240 is not shown in this drawing.) Lines 2232, which are shown as separate (and would be in a diode array), may instead be a continuous sheet, and provide the ground connection.\nFIG. 2G shows an example of resistance (R) over time (t) for a single PCM cell following a single PCM write event at time t=0. The resistance curve 2400 for a cell which has been reset (i.e. which is in its high-resistance state) may rise at first, but then drifts significantly lower. The resistance curve 2410 for a cell in the Set state is much flatter. The sense margin 2420, i.e., the difference between set and reset resistances, also decreases over time. Larger sense margins generally result in more reliable reads, and a sense margin which is too small may not permit reliable reading at all. 2G represents the approximate behavior of one known PCM material; other PCM material compositions may behave differently. For example, other PCM material compositions may display variation of the set resistance over time.\nThe downwards drift of reset resistance may be due to, for example, shrinking size of the amorphous zone of the phase-change material, due to crystal growth; and, in some cells, spontaneous nucleation steepening the drift curve (possibly only slightly) due to introducing further conductive elements into the mushroom-shaped programmable region.\nFIG. 2H shows an example of a processing system 2300. Typically, a processing system 2300 will incorporate at least some of interconnected power supplies 2310, processor units 2320 performing processing functions, memory units 2330 supplying stored data and instructions, and I/O units 2340 controlling communications internally and with external devices 2350.\nFIG. 2I shows an example of a PCM single-ended sensing memory. Two different PCM cells 2400 on different ends of a sense amplifier can be selected separately. Selected elements 2410 are separately sensed by a single-ended sense amplifier 2420.\nFIG. 2J shows an example of a known PCM single-ended sense amplifier 2500. Generally, in a single-ended sense amplifier, a cell read output conducted by a selected bitline BLB is compared against a reference current to provide a digital output OUT. When the PRECHARGE signal turns on transistor 2530, voltage V04 (e.g., 400 mV) precharges the bitline BLB. After precharge ends, the READ signal turns on transistor 2550. Transistor 2550 is connected, through source follower 2560 and load 2580, to provide a voltage which comparator 2600 compares to Voltage_REF, to thereby generate the digital output OUT.\nA variety of nonvolatile memory technologies have been proposed over recent decades, and many of them have required some engineering to provide reference values for sensing. However, the requirements and constraints of phase-change memory are fundamentally different from those of any other kind of nonvolatile memory. Many memory technologies (such as EEPROM, EPROM, MNOS, and flash) test the threshold voltage of the transistor in a selected cell, so referencing must allow for the transistor's behavior. By contrast, phase-change memory simply senses the resistance of the selected cell. This avoids the complexities of providing a reference which will distinguish two (or more) possibilities for an active device's state, but does require detecting a resistance value, and tracking external variations (e.g. temperature and supply voltage) which may affect the instantaneous value of that resistance.\nThe possibility of storing more than one bit of data in a single phase-change material has also been suggested. Phase-change memories implementing such architectures are referred to here as “multibit” PCMs. If the “Set” and/or “Reset” operations can be controlled to produce multiple electrically distinguishable states, then more than one bit of information can be stored in each phase-change material location. It is known that the current over time profile of the Set operation can be controlled to produce electrically distinguishable results, though this can be due to more than one effect. In the simplest implementation, shorter anneals—too short to produce full annealing of the amorphous layer—can be used to produce one or more intermediate states. In some materials, different crystalline phases can also be produced by appropriate selection of the current over time profile. However, what is important for the present application is merely that electrically distinguishable states can be produced.\nFor example, if the complete layer of phase-change material can have four possible I/V characteristics, two bits of information can be stored in each cell—IF the read cycle can accurately distinguish among the four different states.\n(The I/V characteristics of the cells which are not in the fully Set state are typically nonlinear, so it is more accurate to distinguish the states in terms of current flow at a given voltage; resistance is often used as a shorthand term, but implies a linearity which may not be present.)\nIn order to make use of the possible multibit cell structures, it is necessary to reliably distinguish among the possible states. To make this distinction reliably, there must be some margin of safety, despite the change in characteristics which may occur due to history, manufacturing tolerances, and environmental factors. Thus, the read architecture of multibit PCMs is a far more difficult challenge than it is for PCMs with single-bit cells."}
{"text": "Robotic devices (i.e., “robots”) are capable of performing many tasks, such as data collection, searching, and mapping. Generally, these tasks are defined programmatically using a software development system. One or more robots receive the programmed instructions, execute them, and perform the corresponding tasks. This generally occurs under the auspices of a central command and control apparatus.\nWhen a task is complex or large scale (e.g., mapping a big area), the time to complete the task can be prohibitive if an insufficient number of robots are employed. Consequently, to improve efficiency, more robots are typically added to share the task (or parts of the task) and drive it to completion earlier. A small increase in the number of deployed robots can complicate command and control, but not to the point where the array of robots becomes unmanageable.\nDifficulties can ensue when efficient completion of the task requires the deployment of more than just a small number of robots. At this level, the central command and control apparatus typically cannot exercise complete management of the group of robots. Further, a human operator overseeing the robots can be overwhelmed trying to monitor the progress of (and potentially control) every robot. In fact, when working with large numbers of robots, an operator often cannot manually program, charge, or even turn on the robots. Moreover, control strategies for multiple robots need to be robust in the face of complex environments and tolerant to the failure of any number of individual robots within the group. Optimally, a control strategy is designed to be completely scaleable to function with any number of robots.\nA robot typically observes its environment during the performance of the task. For example, a robot can sense when it has encountered an obstacle and, optionally, attempt to circumvent the obstacle. When several robots are deployed to complete a task, it is likely they will encounter each other. On the one hand, the robots can simply avoid each other (i.e., to eliminate the potential of collisions) and continue to perform their duties in a virtually independent manner. Nevertheless, efficiency can be improved beyond that attained by increasing the number of robots by allowing the robots to interact with each other and adapt their behavior accordingly, ideally without intervention from a central command and control apparatus or an attending operator (i.e., autonomously).\nFrom the foregoing, it is apparent that there is still a need for a way to allow increasing the number of robots performing a particular task while, at the same time, distributing command and control to, for example, the individual robots. Further, it is desirable to accomplish this in conjunction with allowing interaction between the robots, resulting in the adaptive behavior thereof."}
{"text": "Many organizations utilize demographic, and geographically-based demographic information, for a wide variety of purposes. For example, the United States conducts a constitutionally mandated, geographically-based census survey every ten years for purposes of readjusting governmental and political representation districts. Demographic information available from the census includes age, gender, ethnicity, race, income, education, etc. The census is conducted within each state plus the District of Columbia, with information provided at a comparatively fine-grained level, with the U.S. geographically divided into over 3,000 counties, which are further geographically subdivided into over 65,000 census tracts, which are further geographically subdivided into over 208,000 census block groups, and which are further geographically subdivided into over 8,200,000 census blocks.\nBased on the notion that demographics are often skewed by geographies, such as wealthy neighborhoods and ethnic neighborhoods, “geo-demographics” have often been utilized to target or preferentially communicate with certain groups, such as for political and marketing communications. In other circumstances, governmental and business decisions are based on such geo-demographics, such as decisions for locating a new post office, a new dealership, a new shopping center or store, and so on. In addition, for certain industries, such as the broadcast and print media, their usage is specifically skewed by such geo-demographic “segments”, as such media may only have applicability for certain demographic groups within certain defined geographies.\nAs a consequence, geo-demographic “clustering” has evolved as a way to try to identify and describe groups of individuals based upon their demographic characteristics. Certain regions are then defined by the corresponding demographic clusters contained in that region. While certainly better than “blind” targeting or analysis by at least providing some differentiation (or discriminatory) capability, such clustering techniques do not provide other significant information, such as behavioral and attitudinal information. (“Discriminatory”, as used herein, should be utilized in a statistical sense, to provide differentiation on the basis of certain common characteristics, such as differentiating people based upon their age groups or educational levels.)\nIn addition, such clustering techniques tend to be either under inclusive or over inclusive. For example, such clustering is under inclusive, as small but significant groups of people in certain regions may be omitted from analysis altogether, as not falling within the major cluster groups for those regions. As a consequence, such clustering systems lose population nuances and minority influences, which may, in fact, be significant. For example, current clustering techniques from a wide variety of companies would group Hispanic individuals of Dominican Republic origin into a cluster residing in large cities; actual analysis would reveal, however, that many of these cities having such a “large city cluster” do not, in fact, have any significant population of such Hispanic individuals of Dominican Republic origin.\nThe over inclusiveness of such clustering groups is also well-known. Groups of individuals with very different behaviors and attitudes, for example, may be grouped together. For example, a young or middle-aged statistician with a doctoral level education, who happens to live in Florida, may, by virtue of his residence, be included within a cluster of retired individuals. Because of such over inclusiveness, nuances and minority information of certain regions and neighborhoods are lost in these forms of cluster analysis.\nAs a consequence, a need remains for other forms of analysis, which not only provide geo-demographic information, but which also provide important behavioral (e.g., what people say and do) and attitudinal (e.g., beliefs, outlooks, psychographics) information. Such an analytical system should have a firm foundation in empirical research, such as through the use of statistically significant or relevant surveys from a nationally represented sample of individuals.\nSuch systems and methods should be fine-grained, and should accurately preserve behavioral nuances and minority information. Such methods and systems should have the capability to model all regions and their subregions, even those that have an insufficient sample size within nationally representative surveys."}
{"text": "Timber mats are commonly used for temporary or portable road ways, bridge decking, excavation support, crane pads and pipeline construction. Timber mats usually are made of planks, which can be oak, mixed hardwoods, Douglas fir, or man-made materials. The timber mats must withstand large forces such as the weight of heavy trucks and machinery positioned or moving on top of them, and so must be supportive and secure. To attach the planks securely together, metal rods or bolts are generally driven through the planks."}
{"text": "This invention relates to a photoelectric tilt-detecting sensor for detecting the presence or absence of tilting of an object to be detected, and a method of fabricating the same.\nThere have been known photo-interrupters for use in various applications of optical sensors. The photo-interrupter generally includes a light-emitting element for emitting light, and a light-receiving element for receiving the light emitted therefrom. The photo-interrupter is utilized to determine whether or not a certain object is present between these two optical elements, by detection of an electric current generated from reception of light by the light-receiving element.\nThe conventionally known photoelectric tilt-detecting sensor, as shown in FIG. 13, comprising a photo-interrupter 31 formed by a light-emitting element 31a and a light-receiving element 31b. The photo-interrupter 31 is built in a lower case 32 for the sensor formed of a resin, with an upper case 33 fitted over the lower case 32. The upper case 33 has opposite bearing supports 33a extending therefrom toward the inner side of the lower case 32. A fan-shaped light-shielding plate 34 is pivotally supported on bearings 33b for swinging between the light-emitting element 31a and the light-receiving element 31b.\nIn the tilt-detecting sensor thus constructed, the light-shielding plate 34 ordinarily rests at its lowest position due to the action of gravity. To use the tilt-detecting sensor, it is mounted or fixed on an object to be detected, not shown. While the tilting angle is smaller than a certain value, the light emitted from the light-emitting element 31a is shielded by the light-shielding plate 34 and no light is transmitted to the light-receiving element 31b. However, if the angle of tilting exceeds the above angle, the light-shielding plate 34 goes out between the light-emitting element 31a and the light-receiving element 31b. As a result, the light emitted is transmitted to the light-receiving element 31b. The light-receiving element 31b, upon reception of light, generates an electric current to supply the same through lead terminals thereof, making possible detection of the tilting for the object to be detected.\nMeanwhile, there is increase in demands of offering sensors integral with a photo-interrupter, as the technology in the art of semiconductor or optical devices advances. On the other hand, it is strongly demanded that the number of parts be reduced for making such devices, together with simplification of fabrication processes.\nHowever, the conventional tilt-detecting sensor necessitates bearings 33b on the upper case 33 for assembling a light-shielding plate 34 therein. Besides the light-shielding plate 34 has to be assembled by inserting its support axis into the respective bearings 33b. This light-shielding plate 34 must be formed relatively complicated in shape, or otherwise assembled with a separate support shaft, thus incurring complexity in fabrication processes. In practical, there is difficulty of reducing the number of parts used and simplifing the fabrication processes.\nFurther, in the conventional tilt-detecting sensor, there is a problem left in smoothness of movement of the light-shielding plate which is in pivotal support by the opposite bearings. The light-shielding plate thus mechanically structured is apt to lose its smoothness in movement, particularly where an external impact force be applied or tilting is in other direction than the direction being detected."}
{"text": "In recent years, as the number of users utilizing smartphones has explosively increased, the demand for mobile Internet services and multimedia has risen sharply. This has given impulse to active research, standardization, and development of improved communication systems. It is necessary for an improved communication system to have a high data rate and high frequency efficiency.\nVarious schemes to achieve high data rates and high frequency efficiency are under active discussion. Examples of such schemes may include: single-user and multi-user MIMO transmission, Coordinated Muli-Point (CoMP) technology aiming to increase the mean data rate within a cell and cell-edge user performance through effective interference control in cell boundary areas, and new cellular network architecture supporting a heterogeneous cell structure composed of a macrocell, picocell and femtocell to prevent performance degradation in a shadow area of an existing macro cell and to guarantee performance at a hot spot where heavy user data is generated.\nIn a multicell scheme raising the level of reuse of frequency and spatial resources in a mobile/wireless communication system, when communication entities transmit data independently of each other, it is inevitable for a user equipment (UE) to receive an interference signal. In particular, as a user equipment moves toward the cell edge region, the interference signal power increases and the data rate toward the base station sharply decreases.\nTo mitigate inter-cell interference, the base station may use a beamforming technique to adjust interference toward other cells, reducing the impact of interference on the system. Here, suitable user equipment selection by the base station may further increase overall system throughput by providing additional gains in addition to beamforming, i.e. further reducing inter-cell interference and increasing signal strength toward a desired user equipment."}
{"text": "Silencer devices for attenuating fluid borne noises are not, in and of themselves, new. Examples of heretofore known silencer devices of this type are shown in U.S. Pat. Nos. 2,940,537 issued on June 14, 1960 to C. D. Smith et al; 2,988,302 issued on June 13, 1961 to C. D. Smith; and 3,019,850 issued on Feb. 6, 1962 to J. J. March. In those heretofore known silencer devices having flow tubes attached to headers, the headers are of unitary construction material. Thus, in order to manufacture a family of silencer devices of different sizes, i.e., having differing numbers of flow tubes from one another, different manufacturing tooling, such as, for example, a metal forming die, must be used for each different sized header. This requires a proliferation of tooling and is therefore costly. Because it is costly and time consuming to set up manufacturing tooling, it is more efficient to manufacture more headers of a given size than may be needed at the time of manufacture and to put the excess number of headers in inventory for future use when needed. This solution is also expensive."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of automatic testing of active plates for liquid crystal displays.\n2. Prior Art\nLCD panels typically are formed with a liquid crystal material sandwiched between an active plate and a ground plate. Polarizers, colorizing filters and spacers may also be included between the plates. During fabrication, many active panels may be formed on a single glass plate. In each area of the glass plate that is to form an active panel, pixel areas, drive lines, gate lines and drive elements are formed. Typically, thin-film transistors are used for the drive elements.\nBecause of the relative complexity of the active plate in comparison to the ground plate, most LCD display defects may be traced to some form of defect in the active plate. Because of the substantial additional expense to fabricate operative LCD displays from the active plates, various techniques have been developed for inspecting the active plates alone so that defective active plates can be identified and repaired or discarded at that stage of the fabrication process.\nAn exemplary portion of an active plate 20 for a monochrome display is schematically illustrated in FIG. 1. As may be seen therein, a plurality of conductive areas 14 are arranged in an XY matrix, each conductive area defining a pixel in the final display. Associated with each conductive area is a thin film transistor 16 having its input terminal connected to a respective row line in the matrix, and its control terminal or gate connected to a respective column line. In the particular matrix shown in FIG. 1, adjacent row lines are coupled to opposite sides of the matrix, and adjacent column lines are coupled to the top and bottom of the matrix, respectively. A color display may be similar, though each pixel of the display is comprised of 3 pixels on the active plate, each for a different color.\nIn operation, each conductive area is charged to a voltage on the respective array lines when the respective thin film transistor is turned on by the voltage on the respective column line during each scan of the array. However, while the active plate may be electrically exercised without the ground plate and the liquid crystal material there between, no visibly perceivable changes occur during that exercise.\nVarious techniques are well known in the prior art for inspecting and testing LCD active plates before proceeding with the further fabrication of the full LCD panel. Each of these techniques is well suited to the detection of certain faults, but not suited to the detection of certain other faults.\nOne inspection technique that is commonly used is inspection under a visible light using a digital camera and computer based image analysis. Because the active plate is comprised of an array of a large number of pixel areas and thin film transistors, one convenient form of optical inspection is to form a difference image between repetitive patterns in the image. This is commonly called Automatic Optical Inspection (AOI). Without defects, the difference image is zero. With defects, the difference image is positive or negative. In this way, various faults may be detected such as shorts and opens and other faults in geometry that would cause failure of the display to operate, or unacceptable variations in the image intensity across the pixel array. However, other potential faults are not detectable in this manner, such as some faults which cause one or more transistors to not operate properly and/or which prevent pixel conductive areas from maintaining a charge until refreshed on the next scan.\nVisible light systems may be either transmissive or pass through (the illumination and camera being on opposite sides of the plate), or reflective (camera and lighting on the same side of the plate), with the active plate fixed and the camera on a transport system for step and repeat operation, particularly with a lighting system moving with the camera to assure uniform lighting for each camera view. Normally the camera is also provided with the Z axis transport system for greater flexibility.\nOther known systems for evaluating active plates at that stage of the LCD fabrication process comprise methods of electrically testing the active plate to enable accept/reject decisions based on reasonably accurate projections on how the plate will perform in the finished LCD display. One such technique utilizes a voltage imaging sensor, such as that described in U.S. Pat. No. 4,983,911, assigned to Photon Dynamics, Inc., assignee of the present invention. These systems provide a two-dimensional image of the voltage distribution across a surface of the active plate, thereby allowing the image to be digitized by an appropriate camera. The systems illuminate the active plate with a collimated beam of optical energy of known polarization through an electro-optic modulator disposed sufficiently close to the plate so as to be affected by the voltages on the active plate. Thus the voltage imaging sensor simulates the top half of the LCD panel and converts charge (voltage) on the active plate to visible light. Accordingly, such a system requires at least the electro-optic modulator to be positioned very close to the substrate. This approach can detect faults in operation, such as defective transistors, but has a resolution limit higher than an AOI system.\nSystems of the voltage imaging type are manufactured by Photon Dynamics, Inc., assignee of the present invention. Such systems may be plate scale in size or smaller, normally with a camera operating in a step and repeat mode. Both reflective and pass-through systems are known, though normally the reflective mode is used.\nAnother well known testing technique is the electron beam scanning or e-beam scanning technique. In these testing systems, the active plate is placed in a vacuum chamber and the plate scanned with an electron beam, with the secondary electrons being detected by a scintillation camera. The camera output versus electron-beam position provides the image data for the plate voltage. Typically, a small part of the active plate is scanned at any one time, with the entire plate being scanned in a step and repeat process. Insufficient secondary electrons are indicative of a plate defect.\nFinally, active plate testing systems based on charge sensing are also known. These systems are based on the concept of turning on each transistor on the active plate to charge the respective pixel conductive area to a specific voltage, then turning the transistor off, and then turning the transistor on again to short the pixel conductive area to ground while measuring the charge returned from the conductive pixel area. An insufficient charge is indicative of a plate defect. Typically, the transistor on periods while charging the pixel conductive area is on the order of the time the transistor would be on for such purpose in the completed display, and the time before shorting out the pixel conductive area for charge measurement is on the order of the time between scans in the completed LCD.\nThus, in an AOI system using visible light, the camera is typically spaced substantially away from the active plate being inspected, whereas in the voltage imaging sensor type systems, at least the electro-optic modulator must be positioned very close to the active plate to obtain a voltage image of a useful resolution. In an electron beam system, a vacuum environment must be provided, though a camera for providing a two-dimensional image is not required. Instead, only a light sensor is needed, preferably with a photomultiplier to increase the light intensity, with the XY information for a two-dimensional image being provided by the electron beam scan control system. Finally, in the charge sensing type equipment, no camera at all is used, and no mechanical transport system is required beyond that used for loading and unloading the active plate to be tested, such as generally also needed in some form for all of the other testing and inspection systems."}
{"text": "This invention relates to stands for supporting various items of clothing and/or recreational equipment. There are many kinds of stands used for this purpose, from a simple post having hooks to fairly elaborate hall trees.\nIn connection with the game of ice or street hockey, there are clothing items, such as helmets, jerseys and pants, pads of various kinds, and the hockey skates or roller blades, all of which must, at times, be temporarily stored. Anyone having a child or teenaged hockey player in the house is well aware of the clutter such items can create if not properly stored.\nIt is an object of the present invention to provide a unique and attractive stand which can be used to support the various items of hockey equipment and clothing and which can provide some motivation for the hockey player to hang up his or her various items of clothing and equipment.\nIt is another object of the present invention to provide such a stand, which is readily portable and which can be easily stored when not in use."}
{"text": "1. Technical Field\nThe present invention relates generally to systems and methods for providing accelerated loading of operating system and application programs upon system boot or application launch and, more particularly, to data storage controllers employing lossless and/or lossy data compression and decompression to provide accelerated loading of operating systems and application programs.\n2. Description of the Related Art\nModern computers utilize a hierarchy of memory devices. To achieve maximum performance levels, modern processors utilize onboard memory and on board cache to obtain high bandwidth access to both program and data. Limitations in process technologies currently prohibit placing a sufficient quantity of onboard memory for most applications. Thus, in order to offer sufficient memory for the operating system(s), application programs, and user data, computers often use various forms of popular off-processor high speed memory including static random access memory (SRAM), synchronous dynamic random access memory (SDRAM), synchronous burst static ram (SBSRAM). Due to the prohibitive cost of the high-speed random access memory, coupled with their power volatility, a third lower level of the hierarchy exists for non-volatile mass storage devices.\nFurthermore, mass storage devices offer increased capacity and fairly economical data storage. Mass storage devices (such as a “hard disk”) typically store the operating system of a computer system, as well as applications and data and rapid access to such data is critical to system performance. The data storage and retrieval bandwidth of mass storage devices, however, is typically much less as compared with the bandwidth of other elements of a computing system. Indeed, over the last decade, although computer processor performance has improved by at least a factor of 50, magnetic disk storage performance has only improved by a factor of 5. Consequently, memory storage devices severely limit the performance of consumer, entertainment, office, workstation, servers, and mainframe computers for all disk and memory intensive operations.\nThe ubiquitous Internet combined with new multimedia applications has put tremendous emphasis on storage volumetric density, storage mass density, storewidth, and power consumption. Specifically, storage density is limited by the number of bits that are encoded in a mass storage device per unit volume. Similarly mass density is defined as storage bits per unit mass. Storewidth is the data rate at which the data may be accessed. There are various ways of categorizing storewidth in terms, several of the more prevalent metrics include sustained continuous storewidth, burst storewidth, and random access storewidth, all typically measured in megabytes/sec. Power consumption is canonically defined in terms of power consumption per bit and may be specified under a number of operating modes including active (while data is being accessed and transmitted) and standby mode. Hence one fairly obvious limitation within the current art is the need for even more volume, mass, and power efficient data storage.\nMagnetic disk mass storage devices currently employed in a variety of home, business, and scientific computing applications suffer from significant seek-time access delays along with profound read/write data rate limitations. Currently the fastest available disk drives support only a sustained output data rate in the tens of megabytes per second data rate (MB/sec). This is in stark contrast to the modern Personal Computer's Peripheral Component Interconnect (PCI) Bus's low end 32 bit/33 Mhz input/output capability of 264 MB/sec and the PC's internal local bus capability of 800 MB/sec.\nAnother problem within the current art is that emergent high performance disk interface standards such as the Small Computer Systems Interface (SCSI-3), Fibre Channel, AT Attachment UltraDMA/66/100, Serial Storage Architecture, and Universal Serial Bus offer only higher data transfer rates through intermediate data buffering in random access memory. These interconnect strategies do not address the fundamental problem that all modern magnetic disk storage devices for the personal computer marketplace are still limited by the same typical physical media restrictions. In practice, faster disk access data rates are only achieved by the high cost solution of simultaneously accessing multiple disk drives with a technique known within the art as data striping and redundant array of independent disks (RAID).\nRAID systems often afford the user the benefit of increased data bandwidth for data storage and retrieval. By simultaneously accessing two or more disk drives, data bandwidth may be increased at a maximum rate that is linear and directly proportional to the number of disks employed. Thus another problem with modern data storage systems utilizing RAID systems is that a linear increase in data bandwidth requires a proportional number of added disk storage devices.\nAnother problem with most modern mass storage devices is their inherent unreliability. Many modern mass storage devices utilize rotating assemblies and other types of electromechanical components that possess failure rates one or more orders of magnitude higher than equivalent solid-state devices. RAID systems employ data redundancy distributed across multiple disks to enhance data storage and retrieval reliability. In the simplest case, data may be explicitly repeated on multiple places on a single disk drive, on multiple places on two or more independent disk drives. More complex techniques are also employed that support various trade-offs between data bandwidth and data reliability.\nStandard types of RAID systems currently available include RAID Levels 0, 1, and 5. The configuration selected depends on the goals to be achieved. Specifically data reliability, data validation, data storage/retrieval bandwidth, and cost all play a role in defining the appropriate RAID data storage solution. RAID level 0 entails pure data striping across multiple disk drives. This increases data bandwidth at best linearly with the number of disk drives utilized. Data reliability and validation capability are decreased. A failure of a single drive results in a complete loss of all data. Thus another problem with RAID systems is that low cost improved bandwidth requires a significant decrease in reliability.\nRAID Level 1 utilizes disk mirroring where data is duplicated on an independent disk subsystem. Validation of data amongst the two independent drives is possible if the data is simultaneously accessed on both disks and subsequently compared. This tends to decrease data bandwidth from even that of a single comparable disk drive. In systems that offer hot swap capability, the failed drive is removed and a replacement drive is inserted. The data on the failed drive is then copied in the background while the entire system continues to operate in a performance degraded but fully operational mode. Once the data rebuild is complete, normal operation resumes. Hence, another problem with RAID systems is the high cost of increased reliability and associated decrease in performance.\nRAID Level 5 employs disk data striping and parity error detection to increase both data bandwidth and reliability simultaneously. A minimum of three disk drives is required for this technique. In the event of a single disk drive failure, that drive may be rebuilt from parity and other data encoded on disk remaining disk drives. In systems that offer hot swap capability, the failed drive is removed and a replacement drive is inserted. The data on the failed drive is then rebuilt in the background while the entire system continues to operate in a performance degraded but fully operational mode. Once the data rebuild is complete, normal operation resumes.\nThus another problem with redundant modern mass storage devices is the degradation of data bandwidth when a storage device fails. Additional problems with bandwidth limitations and reliability similarly occur within the art by all other forms of sequential, pseudo-random, and random access mass storage devices. These and other limitations within the current art are addressed by the present invention."}
{"text": "In the light of international efforts for reduce pollutant emissions from heating systems, particularly gas turbines, burners and methods of operation for burners have been developed in recent years which have particularly low nitrous oxide (NOx) emissions. In such cases emphasis has frequently been placed on the fact that such burners are able to be operated not just with one fuel but where possible with a wide variety of fuels, for example oil, natural gas and/or synthetic gas (e.g. coal gas) as required or even in combination in order to increase security of supply and flexibility during operation. Such burners are described for example in EP 0 276 696 B1.\nAn associated problem is that of a stable combustion in the burner, which is based on a stable location of the combustion zone in the burner. This alters drastically if changes occur in the composition of the fuel, meaning for example a combustion gas having a high proportion of high-grade saturated hydrocarbons, such as C2+ alkane, ethane or propane. Under such circumstances there is a danger of flame blowbacks in the burner. Patent WO 03/062618 A1 thus in particular monitors the C2+ alkane of the inflow combustion gas through infrared absorption. To avoid a blowback, with an increased proportion of C2+ alkane, the combustion gas characteristic is regulated by intervening for example in the combustion gas supply, but also by injecting water or steam.\nThe problem in designing burners for all possible different operating conditions and fuels, especially when there is also a variation in the fuel composition or if there are fluctuations in fuel quality, lies in the fact that the volumes needed during operation in each case (fuel mass flow) of the individual fuels are completely different, making it difficult to use the same supply system and the same spray openings for all fuels. The use of different supply systems for liquid and gaseous materials is thus known in the prior art. A further problem then presents itself however if alternate gaseous fuels with entirely different specific calorific values, for example natural gas and coal gas, are to be used. The completely different volume ratios when using these two fuels and the different chemical processes during their combustion demand a modification and expansion of the known systems.\nModern low-NOx combustion systems are usually based on the so-called “jet in crossflow” mixing-in concept. Low-pollutant combustion, especially with low NOx emissions, can in this case be undertaken by corresponding designs of the mixing-in of the fuel into the cross-flow combustion air. An important design variable in such cases is the penetration depth of the fuel jets into the cross-flow air. This mechanical design adapted in the best possible manner is then correspondingly only used for a specific fuel composition. The invention now uses as its starting point the problem that, with a temporary alteration to the fuel composition or if the fuel is changed, the result can be an alteration of the mixing field, which, with an unfavorable mixture, usually leads to increased NOx emissions.\nBased on this observation, an object of the invention is to specify a method for operation of a burner with which a low nitrous oxide combustion is possible even if there is a change to the fuel composition. A further object of the invention is to specify a suitable device for executing the method.\nThe first object is achieved in accordance with the invention by a method for operation of a burner, in which a fuel is supplied to the burner, sprayed into the combustion air and mixed with the combustion air into a fuel/air mixture and burned in a combustion chamber, in which case, to reduce the nitrous oxide emissions, a fuel characteristic is set explicitly to a predetermined nitrous oxide emission, with a change in a parameter characterizing the fuel being determined and with a penetration depth of the fuel jets into the combustion air adapted to the change being effected.\nIn this case the invention starts from the knowledge that the influence of fuel composition fluctuations on the NOx emissions should if possible not be compensated for by expensive burner-side measures or by expensive adaptations to the mechanical design in the combustion chamber, in order to achieve sufficient flexibility and timely adaptation to the predetermined nitrous oxide emissions. Constructive measures could only reduce the sensitivity to fluctuations in fuel quality—and thereby mixing field variations—to a limited extent, but could not eliminate them completely. This is attributable to the fact that, for a quite specific fuel composition in each case (fuel characteristic) the injection and mixing-in facilities of the burner are “optimized”. Designs with a plurality of injection spray points for the fuel—well distributed over the cross section through which the flow is conducted—or the use of static mixers for setting a desired mixture field, in the absence of the required flexibility—especially with shorter-term variations in fuel composition—are not suitable for guaranteeing compliance with permitted emission limits for NOx emission during operation of the burner or of the combustion system. On the other hand the invention starts from the knowledge that, by spraying the fuel into the combustion air with a most favorable possible penetration level of the fuel jets, the mixing field will be adjusted in respect of a low-emission combustion. This mixing field can even be maintained during operation, taking into account the parameter characterizing the fuel.\nThus the invention proposes for the first time achieving an especially low nitrous oxide combustion by monitoring the fuel characteristic of the fuel composition, in order, if required, to use suitable measures to once again set an optimum low-pollutant operating mode as regards nitrous oxide if a parameter characterizing the fuel changes. This creates the opportunity of keeping the nitrous oxide emissions below a predetermined limit using a parameter which characterizes the fuel, with a penetration depth of the fuel jets adapted to the change being effected in the combustion air.\nPreferably in this case the fuel is sprayed into the combustion air and mixed with the combustion air. Fuel and combustion air are mixed in the burner, with the best possible penetration depth of the fuel rays into the combustion air having to be ensured for the injection of the fuel into the combustion air. This means that the mixing field can be adjusted in respect of a low-pollutant combustion and can be maintained even during operation taking into account the parameter characterizing the fuel.\nIn an especially preferred embodiment a change in the parameter characterizing the fuel is registered and transferred to a control system. In this case the parameter characterizing the fuel is preferably continuously detected and evaluated in the control system. The parameter characterizing the fuel can in this case be determined by a suitable measurement of the fuel flow during operation and in this way the timing of the parameter characterizing the fuel can be stored and evaluated.\nPreferably the fuel characteristic is set explicitly, with the value being set to a reference or required value of the parameter characterizing the fuel at which the predetermined pollutant emission occurs. In this case characteristic performance data determined in advance can already be stored in the control system, with said data representing the relationship between the fuel composition and the nitrous oxide emissions. Alternatively however an in-situ measurement of both the current nitrous oxide emission values and also the fuel composition is simultaneously measured and transferred to the control system.\nIn an especially preferred embodiment the Wobbe index is determined from a parameter characterizing the fuel. What is referred to as the Wobbe index is a normal standard used to characterize the fuel composition and temperature. The Wobbe index allows a comparison of the heat content of different fuels relative to their volume, especially combustion gases, to be made at different temperatures. Since a combustion system, such as a gas turbine for example, is operated such that in the final analysis heat energy is released in a combustion chamber and that the fuel flow is set by controlling the volume flow, fuels with different fuel composition but still with relatively similar Wobbe indexes are generally supplied by the same fuel supply system to the burner. Variations in the fuel composition lead to variations in the nitrous oxide emissions, with a setting of the Wobbe index enabling a permitted highest nitrous oxide emissions in the operation of the gas turbine.\nIn the inventive method the Wobbe index of the fuel is preferably determined by evaluating\nthe relationship\n            W      ⁢                          ⁢      I        =                  L        ⁢                                  ⁢        H        ⁢                                  ⁢        V                              S          ⁢                                          ⁢                      G            ·                          T              /                              T                Ref                                                          ,with LHV being the lower heat value of the fuel, T its absolute temperature and SG the specific gravity of the fuel relative to the air under standard conditions, and TRef being a reference temperature.\nIn this case, for setting the desired Wobbe index, the temperature of the fuel is preferably set explicitly with respect to a predetermined nitrous oxide emission. According to the above formula the Wobbe index is related in a relatively simple manner to the current fuel temperature, namely inversely proportional to the square root of the fuel temperature. This means that, if the Wobbe index changes, i.e. the Wobbe index deviates from a predetermined required value with low nitrous oxide emissions, the desired Wobbe index and thereby the desired nitrous oxide emission can be set by a corresponding regulation of the fuel temperature. Depending on the situation the fuel can be warmed up or cooled down to the required value to enable a temperature setting of the required value of the Wobbe index to be used to achieve the desired NOx emission.\nIt is however also possible for a medium to be mixed in with the fuel for setting the fuel characteristic. A modification of the Wobbe index is also especially to be achieved in this way in order to ensure a low-pollution operation of the burner. The preferred medium considered for injection into the fuel is water, steam or nitrogen, but also hydrocarbons with a high heat value for example.\nAs an alternative to the Wobbe index, the so-called impulse flow density ratio can also be determined and evaluated as a parameter characterizing the operating state. Impulse flow density ratio: The mixture quality with “jet in crossflow” depends, with a given geometry, on the impulse flow density ratio, i.e. on the quotient of the impulse flow density ratio of the jet and on the impulse flow density ratio of the crossflow.Impulse I=M·c=p·c·A·c Impulse density I=I/A (surface)=pc2 \nImpulse flow density of the air is essentially given by ambient conditions and gas turbine performance. Impulse flow density of the fuel, apart from depending on the gas turbine performance also only depends on the fuel composition. The heating value gives the mass flow and thereby with the density for fixed geometry the impulse flow density. The impulse flow density is thus not a fuel characteristic but a variable which depends on the fuel composition. This variable can also be regulated to a desired required value if it changes by monitoring the fuel composition during the operation of the burner.\nPreferably the method is applied when a burner of a gas turbine is operated. The demands for low-pollutant combustion during operation of gas turbines, especially with stationary gas turbines for energy production, have increased continuously in recent years. The method In accordance with the invention makes low-emission operation possible, with the regulation measures described above already being able to be undertaken on the fuel side during the operation of the gas turbine system if there are variations in the composition of the fuel. Expensive constructive measures on the burner side can be dispensed with.\nA liquid or gaseous fuel is preferably used for operation. The method can be used example with oil, natural gas or with a synthetic gas, e.g. coal gas.\nAn object of the invention directed to a device is achieved by a device for carrying out the method with an analysis device for analysis of the current fuel composition during burner operation and with a checking and control system for determining a deviation and for setting the parameter characterizing a fuel to a required value, at which the predetermined pollutant emission is present.\nThe advantages of the inventive device are produced in the same way as the above-described advantages relating to the method.\nThe checking and control system is preferably designed in this case for setting the fuel temperature of the fuel, i.e. a heating up or cooling down the fuel as required.\nThe checking and control system also preferably includes means for controlled injection of an inert medium, especially steam, water or nitrogen or of a hydrocarbon, into the fuel.\nThe same reference symbols have the same meaning in the figures."}
{"text": "1. Field of the Invention\nThe field of the invention relates to microelectromechanical systems (MEMS).\n2. Description of the Related Technology\nMicroelectromechanical systems (MEMS) include micro mechanical elements, actuators, and electronics. Micromechanical elements may be created using deposition, etching, and or other micromachining processes that etch away parts of substrates and/or deposited material layers or that add layers to form electrical and electromechanical devices. One type of MEMS device is called an interferometric modulator. As used herein, the term interferometric modulator or interferometric light modulator refers to a device that selectively absorbs and/or reflects light using the principles of optical interference. In certain embodiments, an interferometric modulator may comprise a pair of conductive plates, one or both of which may be transparent and/or reflective in whole or part and capable of relative motion upon application of an appropriate electrical signal. In a particular embodiment, one plate may comprise a stationary layer deposited on a substrate and the other plate may comprise a metallic membrane separated from the stationary layer by an air gap. As described herein in more detail, the position of one plate in relation to another can change the optical interference of light incident on the interferometric modulator. Such devices have a wide range of applications, and it would be beneficial in the art to utilize and/or modify the characteristics of these types of devices so that their features can be exploited in improving existing products and creating new products that have not yet been developed."}
{"text": "The present invention relates to a method for preparing a new and improved fire-resistant, insulated conductor and to new and improved insulated conductors made in accordance therewith. More particularly, it relates to a fire-resistant insulation prepared by helically wrapping a conductor with at least one layer of mica-impregnated glass tape in a half-lapped fashion to form a continuous first layer, and applying a continuous outer layer of a curable, platinum catalyzed silicone rubber composition thereto and curing said outer layer to form a hardened, fire-resistant insulation.\nIn recent times, the hazards associated with fire in skyscrapers and other heavily populated buildings have become more appreciated. The need to provide circuit integrity to certain critical circuits, such as fire alarms and elevators, in the event of fire cannot be over-emphasized.\nIt has long been known to use mica in insulating compositions. Having excellent dielectric properties and fire-resistance, this natural material is well suited for use in electrical insulation applications. Typically, mica insulations are provided in the form of mica tapes, such as described in U.S. Pat. No. 2,656,290. As described therein, individual mica flakes are bonded to one another, as well as to a pliable base sheet, and, if desired, also a cover sheet, by a liquid bonding agent which may be hardened by suitable additives. The bonded mica tape used for these purposes may be relatively narrow, having a width of 2 to 3 cm for example, or it may be used in sheets of greater width. A selected conductor is wrapped with the mica tape and the wrapped conductor is subjected to a vacuum and impregnated with a thin liquid impregnating resin. The resin and the bonding agent are specifically chosen such that the bonding agent, together with the hardners and polymerization accelerators present in the impregnating resin, combine completely with the impregnating resin to form a uniform hardened insulative coating. Considerable time and effort has been spent in making mica tape insulation which is flexible and exhibits good adhesion of mica to tape. The shortcoming with such tapes is that the vacuum impregnation step tends to be costly and care must be taken that the impregnating resin is fully dispersed throughout the windings to eliminate voids in the insulation which decrease the dielectric properties of the resulting insulation. Further, these compositions generally are not fire-resistant enough at high temperatures to be suitable in today's applications.\nMore recently it has been known to use fluoropolymers, such as tetrafluoroethylene (TEFLON) and ethylenetetrafluoroethylene (TEFZEL) as insulative coatings. Such materials are effectively fire-resistant up to about 250.degree. C.-300.degree. C. which is an improvement, and have excellent dielectric strength. A serious shortcoming with these insulations however is that at high temperatures when these polymers burn, they may give off toxic gases rendering their use less than desirable in today's construction applications.\nRecently, elsewhere in the polymer art, there have been developed new and improved silicone elastomers such as those described in U.S. Pat. No. 4,061,609, assigned to the same assignee as the present invention. The silicone elastomers described therein include a vinyl-containing polysiloxane, a hydrogen-containing siloxane, a platinum catalyst and further include an inhibitor compound containing at least one hydroperoxy radical added to improve processability. Similar but uninhibited silicone elastomers having a silica filler have been shown to be flame-retardant in U.S. Pat. No. 3,514,424.\nIt has now been found that the platinum-catalyzed silicone elastomers described in the former patent are well suited to produce new and improved fire-resistant insulations in combination with a mica-impregnated glass tape."}
{"text": "1. Field of the Invention\nThe present invention relates to a fuel supply apparatus for an internal combustion engine, and more particularly relates to a fuel supply system including a fuel filter.\n2. Description of Related Art\nIn a conventional fuel supply system for an internal combustion engine, as illustrated in FIG. 8 (as disclosed in Japanese Patent Publication HEI 6-129325), a fuel pressure regulator 108 is disposed in a fuel pipe 106 connecting a fuel pump 102 and a fuel injection valve 104. A high-pressure fuel filter 116 is disposed in a portion of the fuel pipe 106 between the fuel pressure regulator 108 and the fuel injection valve 104. Excess fuel is returned from the fuel pressure regulator 108 to a fuel tank 114 via a fuel-return pipe 110. A low-pressure suction filter 112 is provided to a suction port of the fuel pump 102.\nHowever, the conventional fuel supply system has the following problems:\nFirst, because the high-pressure fuel filter is disposed in a high-pressure portion of the fuel pipe between the fuel pressure regulator and the fuel injection valve, sealing requirements at pipe connecting portions are severe (leakage at the pipe connecting portions is not allowed because fuel leakage will seriously affect fuel injection), so the cost of piping is increased.\nSecond, because the high-pressure fuel filter is disposed in the portion of the fuel pipe between the fuel pressure regulator and the fuel injection valve, blockage of the high-pressure fuel filter and a change in pressure drop at the high-pressure fuel filter due to such blockage will directly affect fuel injection.\nThird, because a large portion of fuel discharged from the fuel pump is returned to the fuel tank via the fuel-return pipe, and because no fuel filter is provided in the fuel-return pipe, fuel in the fuel tank is not cleaned. As a result, the suction filter is likely to cause blockage. If a coarse filter having large pores is used in the suction filter to prevent the filter from being choked, foreign particles in the fuel are not filtered, thereby advancing abrasion of the fuel pump."}
{"text": "1. Technical Field to Which the Invention Belongs\nThe present invention relates to a vehicle control system for controlling a vehicle on the basis of a road shape, as decided on the basis of map data comprised of a set of coordinate points obtained from a navigation system.\n2. Discussion of Relevant Art\nAccording to a previously proposed system (i.e., Japanese Patent Application Laid-Open No. 6-281471) of the same Assignee, a passable zone is set on a map on the basis of map data obtained from a navigation system, so that when a curve ahead of the subject vehicle of the driver goes out of the passable zone, a warning is issued to urge the driver to decelerate the vehicle, or an automatic braking or the like is actuated to effect an automatic deceleration.\nIn the aforementioned previously proposed system, however, it is decided whether or not the vehicle can pass through a curve far ahead of its actual position. As a result, arithmetic operations, not necessarily indispensable, are performed to increase loads on the arithmetic units. This raises a problem that the arithmetic units must have large capacities. When the road ahead is bent in S-shape, moreover, this S-shaped road may partially belong to an automatic deceleration zone or a warning zone, even if the curve is actually passable because a sufficient deceleration distance is present from the actual position to the curve. Then, the not-necessarily-indispensable automatic deceleration or warning may be executed to bother the driver."}
{"text": "The invention relates to a method for evaluating a substance, especially a hydrocarbon fuel such as gasoline, diesel fuel, naphtha, jet fuel, kerosene and the like, to determine one or more parameters of the substance.\nAccurate methods for measuring parameters of fuels such as gasoline are desirable. It is especially useful to obtain accurate measurements of parameters such as octane numbers, Reid vapor pressure and the like so that products can be supplied to customers which accurately meet the customers specifications.\nNear infrared (NIR) spectroscopy has been found to be useful in evaluating certain substances because certain bands of the NIR signal are indicative of parameters of the substance. For example, U.S. Pat. No. 4,963,745 to Maggard discloses a process for measuring octane of gasoline using NIR spectroscopy and mathematical processing using derivatives. This method, sometimes referred to as multi-variate analysis, is useful in measuring some parameters. However, the method requires a high signal-to-noise ratio and long-term stability, and further requires a linear correlation between the NIR spectrum data obtained and the parameter to be determined.\nIn accordance with the foregoing, it is the principal object of the present invention to provide a method for evaluating desired parameters of a substance, particularly a hydrocarbon such as gasoline or other types of fuel using NIR spectroscopy wherein accurate readings are obtained despite non-linearity of the relation between the input NIR signal and the output parameter to be determined.\nIt is a further object of the invention to provide a method for evaluation which provides accurate determinations despite potentially large amounts of \"noise\" or corrupted data in the NIR signal.\nIt is still another object of the invention to provide a method for evaluation wherein measurements are obtained \"on-line.\"\nOther objects and advantages will appear herein below."}
{"text": "1. Field of the Invention\nThe present disclosure relates to an ink sheet cartridge, and more particularly, to an ink sheet cartridge used in a line thermal printer for printing a print image.\n2. Description of the Related Art\nIn general, a line thermal printer having a longitudinal thermal head, and heater elements of which are arranged in a width direction of a printing paper, is used as a printer which promptly performs a printing operation line-by-line or page-by-page on a printing paper.\nAn ink sheet cartridge which receives an ink sheet unit in a cartridge case is fitted into the line thermal printer as an ink sheet cartridge for the line thermal printer.\nThe ink sheet cartridge includes a longitudinal ink sheet of which a coloring region is continuously applied with ink of colors yellow, magenta, cyan, or the like. One end of the ink sheet is wound around a sheet supply bobbin of a cylindrical shape while the other end of the ink sheet is wound around a winding bobbin of a cylindrical shape. In addition, the cartridge case includes upper and lower halves made of a synthetic resin, which are bonded or engaged with each other by an ultrasonic welding or a snap-in method.\nA known ink sheet cartridge allows the ink sheet unit to receive the sheet supply and winding bobbins to be rotatable in the cartridge case.\nHowever, in an ink sheet cartridge formed by bonding the upper and lower halves with each other so as to form the cartridge case, it is difficult to reuse the cartridge case as the cartridge case needs, inevitably, to be broken to replace a used ink sheet unit with an unused ink sheet unit. Likewise, in an ink sheet cartridge formed by engaging the upper and lower halves with each other in terms of the snap-in method, it is difficult for a common user to reuse the cartridge case as a tool for detaching the engaged halves by the snap-in method, or a certain degree of skill, is necessary to replace the ink sheet unit.\nAdditionally, in the ink sheet unit of the known ink sheet cartridge, the coloring region applied with the ink of yellow, magenta, cyan, or the like color is exposed. Because of this, even when it is possible to detach the upper and lower halves to replace the ink sheet unit, there is a concern of touching a wound part of the ink sheet in the sheet supply bobbin or the sheet winding bobbin with a hand or finger, etc., so that a foreign substance such as dust can be attached to, or a flaw created on, the coloring region of the ink sheet. Therefore, when an image is printed by using the above-mentioned ink sheet, a failure or glitch in the transfer of the ink results in an unsatisfactory image print.\nExamples of such an ink sheet cartridge are disclosed in Japanese Patent No. 3496267 and Japanese Unexamined Patent Application Publication No. 5-155122."}
{"text": "In recent years, a thin film transistor (TFT) including an oxide semiconductor film such as InGaZnO, which is oxide of indium (In), gallium (Ga) and zinc (Zn) has been developed. It is known that the InGaZnO thin film transistor can be easily fabricated at a low temperature, and the InGaZnO thin film has a high mobility of 10 cm2/Vs or more. Therefore, realization of a back end of line (BEOL) transistor, which is embedded in multilayer interconnection, has been expected. Furthermore, application to liquid crystal panels, and the like has been expected. In addition, it is known that the channel region of an InGaZnO-TFT is not an inversion layer but a storage layer unlike a silicon (Si) metal oxide semiconductor field effect transistor (MOSFET) or a Si-TFT.\nMeanwhile, the channel length is getting shorter along with recent fining of patterns, and the shortening of the channel length causes a problem of significant deterioration of the threshold voltage of a TFT. For example, a threshold voltage around 0 V may be changed to −10 V or lower."}
{"text": "As dynamic random access memory (DRAM) technologies are scaled to smaller dimensions, reliability issues arise that require mitigation by careful design. One such issue is the noise or disturbance injected into neighboring cells when a row is activated and refreshed repeatedly. The voltage swings and current driven when a row (called the aggressor) is activated or refreshed can inject noise into neighboring cells (victims), potentially draining a little bit of charge from the victim's storage capacitors. The more often the aggressor row is activated before its victim neighbors are activated or refreshed, and the weaker the victim cells (smaller charge, higher leakage, etc.), the higher the probability that one or more cells in the victim row will flip and thus may cause data corruption. The multiple activates to an aggressor row is typically referred to as row hammer."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of separating (or peeling) a lamination body and a method of manufacturing a semiconductor device formed using a plastic substrate.\n2. Description of the Related Art\nIn recent years, a technique of forming a thin film transistor (TFT) using a semiconductor thin film (with a thickness of from approximately several nm to several hundreds nm), which is formed over a substrate with an insulated surface, has been attracting attention. The thin film transistor has been widely applied in various electronic devices such as an IC and a display device. In particular, development related to the thin film transistor as a switching element for an image display device has been carried out hurriedly.\nVarious applications of such an image display device have been expected, and particularly, application to a portable device has been attracting much attention. A glass substrate and a quartz substrate has been typically used for forming the image display device now, however, these substrates have some drawbacks of being fragile and heavy. Further, these substrates are unsuitable for mass-production since the surface area thereof is difficulty enlarged. Therefore, it has been tried to form a semiconductor element, e.g., a TFT on a substrate having flexibility as typified by a flexible plastic film.\nIn the case of using the flexible plastic substrate, however, the maximum temperature of the process should be lowered since the plastic film has low heat resistant properties. Accordingly, it has been impossible to form a semiconductor element, e.g., a TFT having as good electric characteristics as that formed over a glass substrate. Thus, a high-performance semiconductor device, e.g., a liquid crystal display device or light emitting element using a plastic film has been developed.\nVarious kinds of methods for separating a lamination body, which is formed over a substrate through a separation body, from the substrate have been already proposed. For example, there is a technique as disclosed in patent document 1 and patent document 2, wherein a separation layer is formed on a transparent substrate by using amorphous silicon (or polysilicon), a lamination body is formed thereon, and laser light is irradiated from a side of the substrate to discharge hydrogen contained in the amorphous silicon so that a gap is caused between the separation layer and the substrate and the substrate is separated from the lamination body.\n[Patent Document 1]: Japanese Patent Application Laid-Open No. Hei 10-125929\n[Patent Document 2]: Japanese Patent Application Laid-Open No. Hei 10-125931\nIn the above-mentioned separation method, however, a substrate having a high light-transmitting property is absolutely required. Further, a step for irradiating relatively high-energy laser beam on an entire surface of the substrate is required to apply sufficient energy such that laser beam transmits through the substrate and hydrogen contained in the amorphous silicon is discharged. This might damage the lamination body. In the case where an element is formed on the separation layer according to the above-mentioned separation method, when the element is formed by a heat treatment at high processing temperatures, hydrogen contained in the separation layer is dispersed and reduced. Accordingly, poor separation might be caused even when laser beam is irradiated to the separation layer, which results in reduced yield."}
{"text": "The present invention refers to a clothes washing machine, in particular a household-type clothes washing machine, provided with special means and related operating modes to enable the particular type of textiles, or mix thereof, forming the washload in the drum of said washing machine to be appropriately identified.\nAlthough the present invention relates particularly to front-loading clothes washing machines, and for reasons of greater simplicity and convenience the following description refers to such a type of washing machines, it will be appreciated that the invention may similarly apply to other types of washing machines, such as top-loading washing machines.\nWashing machines are known in the art that are provided with means adapted to identify the type of textiles, or mix thereof, forming the particular washload being handled in the drum of the washing machine. A purpose of such identification is to provide the machine with the ability of selecting the washing cycle automatically, with the various process parameters selected so as to optimize the operation of the machine and the washing results. For instance, the U.S. Pat. No. 5,161,393 to the name of General Electric Company discloses a quite effective method for identifying the type of textiles in the washload. However, such a method only applies to washing machines having their drum rotating about a vertical axis, so that it is not suitable for use in conjunction with the great majority of washing machines having their drum rotating about a horizontal axis, that is, nearly all of the European-built machines. Furthermore, such a method is a sort of a trial-and-error one based on a set of successive measurements, so that it is quite complex and time-consuming."}
{"text": "Currently the major process by which hydrogen is produced is by the steam reforming of methane. Another means of making hydrogen is by the electrolysis of water. The electricity required for electrolysis is mainly derived from the electric power grid, and the predominant source of grid electricity, combustion of fossil fuels, generates emissions, such as nitrogen oxides and particulate matter, as well as carbon dioxide. One way to eliminate such emissions is to use solar generated electricity to electrolyze water to make hydrogen. Presently, efforts are directed toward improving the efficiency, durability, and cost of the solar powered hydrogen production processes.\nHowever, systems consisting of solar cells to make electricity together with electrolyzers to dissociate water into hydrogen and oxygen, as they exist today, cannot produce hydrogen as cheaply as the steam reforming of methane. Several projects have attempted to produce hydrogen gas to supply vehicle-fueling stations by using electricity from photovoltaic panels and commercially available electrolyzers to split water. These projects proved unsatisfactory and were usually short-lived due to the low efficiency and high cost of the combined technology, which only converted about 2%-6% of the solar energy to hydrogen fuel energy, thus, greatly increasing capital costs, the resulting hydrogen fuel cost (at least $11 per kilogram of hydrogen), and the large area covered by the system."}
{"text": "Networked business applications including enterprise resource planning applications (and other business applications) have become prevalent and, as such, software makers have been forced to provide integrated business solutions via on-demand modules and/or on-premise. On-demand refers to locally using a business system which includes one or more modules executed in a cloud (e.g., accessing, via a network/Internet, another physical processor configured to provide the on-demand module(s)). On-premise refers to executing one or more modules of the business system locally. In some instances, the business system may comprise one or more modules (also referred to as components and/or systems) that are implemented on-premises, while other modules may be implemented on-demand. In any case, the use of on-premises and on-demand makes these business systems more valuable because of extended business functionality but at the same time more complex, which drives up the cost and complexity of ownership of the business systems."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to automatic convergence of stereoscopic images based on disparity maps.\n2. Description of the Related Art\nIn human visual systems or stereoscopic camera systems, the point of intersection of the two eye axes or two camera axes is the convergence point. The distance from the convergence point to the eye or camera is the convergence distance. For human eyes, the convergence point can be at an arbitrary distance. For stereoscopic cameras, the convergence point may be, for example, at infinity (for a parallel camera configuration) or at a fixed distance (for a toe-in camera configuration).\nWhen a person looks at stereoscopic images on a stereoscopic display, the eyes naturally converge to the display screen. The distance from the display screen to the eyes is the natural convergence distance. However, to view the 3D effect correctly, the viewer's eyes adjust to have the same convergence distance as the camera. Such constant convergence distance adjustment can cause discomfort over time such as headaches or eye muscle pain."}
{"text": "Microelectronic devices are used in a wide array of products. These devices, such as memory and microprocessor chips and similar devices have traditionally been used in, for example, computers, telephones, sound equipment and other electronic products. Over the last several years, microelectronic devices have become faster, better, and less expensive. Microelectronic devices are accordingly now also used in traditionally non-electronic products, such as appliances, vehicles, toys and games, medical devices, novelty items, etc. The remarkable progress made in the microelectronic device industry has led to improved yet less expensive products of all types. It has also led to entirely new types of products.\nA major factor in the development of microelectronic devices has been the machines and methods used to manufacture them. Manufacturing of microelectronic devices requires extreme precision, extremely pure materials, and an extremely clean manufacturing environment. Even tiny particles of dust, dirt, metals, or manufacturing chemicals, at almost any stage of the manufacturing process, can cause defects and failures in devices. Reducing contaminants is therefore critical to cost effective manufacturing. Accordingly, intensive research and development has focused on reducing contaminants in microelectronic manufacturing processes.\nAs microelectronic devices are made smaller, reducing contaminants has become even more important, and more difficult to achieve. In the past, various approaches have been used to reduce contaminants. These include use of materials that tend not to generate particles, careful selection and placement of mechanical components in processing machines, and use of a flow of highly filtered and clean air or gases, to carry any particles generated away from the wafers or substrates which form the microelectronic devices.\nEven though these techniques have been successful, engineering challenges remain in trying to further reduce contamination, and to provide more reliable and cost effective manufacturing. As described below, the inventors of the new machines and methods described in this patent application have made developments which offer significantly improved manufacturing of microelectronic devices.\nManufacturing of microelectronic devices involves using various chemicals. These chemicals are typically in liquid form, although gases and vapors are also often used. These chemicals must be highly pure and are therefore expensive. Chemicals used in some processes, such as strong acids or oxidizers, are also toxic. Use of these chemicals, and disposal of the chemicals after they are used, can be time consuming and expensive. Consequently, reducing the amount of chemicals used is highly advantageous. On the other hand, in general, enough of the chemicals must be provided so that they can be uniformly applied over all surfaces of the wafer or substrate being processed or manufactured. It can therefore be difficult to minimize chemical consumption, while maintaining good manufacturing results.\nManufacturing microelectronic devices also typically requires large amounts of purified and de-ionized water. After it is used, e.g., for rinsing, the water typically will have small amounts of dissolved chemicals in it. The water then also often requires special handling and disposal efforts. Accordingly, reducing the amount of water used, as well as the amounts of chemicals used, would be highly advantageous.\nOne problem that may arise in existing microelectronic device manufacturing systems is non-uniform processing or etching. Non-uniform etching may occur when a microelectronic workpiece is not precisely aligned within a rotating processing chamber. When a wafer or workpiece is improperly aligned in a rotating processing chamber, non-concentric etching of the workpiece may occur, which can be particularly problematic when the edge of the workpiece is being processed. In this case, a larger edge region is processed on one side of the workpiece diameter than on the other side. As a result, fewer microelectronic devices can be produced from the workpiece. Thus, there is a need for a workpiece processing system that more precisely aligns a workpiece in a processing chamber.\nAnother problem in existing workpiece processing systems arises when processing fluids are not thoroughly removed from the processing chamber, which can result in the mixing of different processing fluids, as well as fluid deposition on the workpiece surfaces. At certain times, even a few drops of excess fluid can result in defects or failure of the microelectronic end products. Thus, there is a need for a workpiece processing system that thoroughly removes processing fluids from the processing chamber."}
{"text": "Field of the Invention\nThe present invention relates to a water output device, and in particular to a temporary stop water output device used in shower equipment.\nThe Prior Arts\nUsually, in using faucet water, sometimes it is required to stop using water temporarily, and then resume using water right after. However, presently, this is realized through a control switch not disposed on the water output equipment itself, and that could cause quite inconvenience.\nTherefore, presently, the design and performance of the water output device is not quite satisfactory, and it leaves much room for improvement."}
{"text": "Bitumen-based membranes are used for waterproofing in construction. Such membranes normally feature a reinforcement made of non-woven polyester cloth or of a layer of reinforced glass such as fiberglass. One of the membrane surfaces is coated with small slate flakes of various colors, similar to granules, for protection against ultraviolet rays.\nThe users of these types of membranes currently face serious problems when trying to seal together contiguous membranes because of the relatively complex steps that must be carried out to connect the two contiguous or adjacent ends. Of course, if an effective seal is not attained, the membrane system will leak and not achieve its very purpose.\nThe purpose of this invention is to solve the above mentioned problem by creating a bitumen-based membrane with sealing means that make it possible to connect continuous membranes quickly and effectively with reduced energy consumption and excellent results.\nWithin the above described scope, one of the goals of the invention is to provide membranes that are impermeably sealed to adjacent or contiguous membranes in an effective and aesthetically pleasing manner.\nThe bitumen-based membrane that is the subject of this invention provides a reliable seal between membranes because of its unique design features.\nThe subject bitumen-based membrane is easy to produce because it is made of elements and materials that are available on the market. Further, the membrane of the present invention is relatively inexpensive to manufacture and install.\nThe task described above, the goals mentioned, as well as additional goals--which will be described hereinafter--have been achieved with a bitumen-based membrane featuring sealing means for contiguous membranes, per this invention, which consists of the following:\nA bitumen-based membrane sheet, developed longitudinally, which is reinforced with a polyester non-woven cloth or with a coat of reinforced glass on one side of said membrane sheet.\nOne side of the membrane sheet will be coated with granular material such as small slate flakes. Said membrane sheet is characterized by the fact that it includes, at least at one of its ends, a transverse selvage which consists of a tract of exposed bitumen, free of granular material, as well as of a thin coat of polyethylene which, when heated, seals the ends of contiguous membranes together."}
{"text": "The present invention relates to an apparatus for sorting packages, of the type comprising means for supplying the packages to a manipulation and orientation station.\nThe present invention further relates to a method for putting the packages in step.\nIn the sector of packaging and palletising plants, a well known need is that of manipulating and/or orienting packages containing products treated on lines for their transfer and/or processing, usually positioned upstream of said manipulating station.\nSaid requirement is particularly acute in plants for treating containers for drinks, for instance cans and bottles, made of glass or plastic, for containing water, beer or wine. In particular, the packages in question may be crates, cartons, bundles or clusters.\nThe patent EP 1046598 discloses an apparatus for manipulating packages of containers, typically bottles or cans, in order to arrange them according to a predefined configuration to allow them to be transported on pallets.\nThis apparatus is provided with a station for aligning and separating the packages, provided with a plurality of parallel and diagonal guides, on each of which slides a grip element. Each grip element takes a package and orients it according to the final configuration to be obtained.\nThe packages are supplied to the aforesaid station in one or more rows, by means of a conveyor belt.\nThe velocity of sliding of the grip elements on the diagonal guides is controlled by a processor and it may vary according to the velocity of advance of the conveyor belt, in order to enable the grip elements to follow the various packages and to manipulate them in such a way as to make them assume a correct positioning angle relative to the direction of motion of the belt.\nThe apparatus described above has the important drawback of requiring as many grip elements are there are rows of packages on the conveyor belt. The grip elements are constrained to move diagonally and therefore they cannot follow packages which simultaneously reach the alignment and separation station. Precisely, said drawback is due to the fact that the grip element have a single actual degree of freedom of motion which corresponds to the diagonal motion along the guide.\nIn accordance with a second known prior art solution, apparatuses exist which are capable exclusively of rotating the packages by means of fixed impediments against which the packages themselves impact asymmetrically. The asymmetry of the impact causes a rotation of each package about its substantially vertical barycentric axis.\nAn apparatus of this type has the important drawback of compromising the integrity of the packages and above all of the objects contained therein. This drawback is particularly felt in the beer processing industry, because the packages contain glass bottles within which is present beer that was just pasteurised (and thus is hot) which, along with the low thickness of the glass, considerably increases the fragility of the bottles and the consequent risk of breakage."}
{"text": "Depending on the mode of operation, hydraulic accumulators are based on the principle of bubble and membrane accumulators or liquid sound dampers. In hydropneumatic dampers, the compressibility of a gas (mostly nitrogen) is used for damping. In the case of a bladder accumulator, for example, the bladder is compressed or expanded as a function of the magnitude of the variations in pressure. A similar behavior is exhibited by membrane accumulators as well. Specifically, hydropneumatic dampers (pulse-tone pulsation dampers) were developed because damping may be impaired because of disadvantageous connection of the hydraulic fluid to the volume of gas when standard bladder or membrane accumulators are used. Dampers of this type have an in-line connecting block by means of which the variations in volume or pressure are optimally coupled to the gas volume. Good damping properties up to a frequency of around 500 Hz can be obtained with these dampers (see Mannesmann-Rexroth design projects “Konstruktion von Hydroanlagen” (Construction of Hydraulic Facilities), Der Hydrauliktrainer, Vol. 3, 1st Edition, page 106).\nIn conventional solutions readily available commercially, a connecting point is provided in the accumulator housing to increase the gas supply volume by mounting on the gas supply side an accumulator cylinder in which the additional operating gas, in particular in the form of nitrogen, is stored. Such solutions are costly in manufacture, and are geometrically large. Sealing problems may arise at the connecting point in question, leading to loss of operating gas. To improve the solutions, it has already been proposed that the gas supply chamber in the accumulator housing itself be enlarged and that additional supply cylinders connected to the accumulator housing be eliminated. However, enlargement of the gas supply chamber in the accumulator housing also increases free, travel for the separating element or separating membrane, so that it is subjected to heavy wear, such as by folding and overstretching processes, so that failure of the accumulator device soon ensues.\nThe solutions, discussed in the foregoing with the accumulator cylinder of operating gas mounted on the device, allow reduction of the clear path of displacement inside the accumulator housing for the separating element. However, on the whole, the operating capacity of the accumulator is adversely affected, especially with respect to pulsation damping.\nIn addition, the conventional solutions have separating elements of an elastomer material which are essentially gas-permeable to a certain extent. Loss of gas due to diffusion and permeation processes occurs in the long term. The gas supply is transferred to the fluid side."}
{"text": "The present invention relates to a hydrogel device containing an active agent which is released by diffusion to an aqueous medium upon contact therewith at a controlled rate for a predetermined period of time, a method for preparing such a device and a method of use of such a device to release the active agent into the aqueous medium.\nVarious devices are known in the art for the sustained release of an active agent. For example, monolithic devices for sustained release of agent, wherein the agent is dispersed uniformly in a non-swellable homogeneous and imperforate polymer matrix where the agent dissolves in and permeates through the polymer itself are known. Microporous devices, in which the pores contain an active agent permeable liquid or gel medium such that the active agent preferentially dissolves in and permeates through the medium in the pores, are also known, as well as osmotic bursting devices wherein water is imbibed osmotically into active agent depots in a water permeable-active agent impermeable polymer matrix to rupture the depots serially. Characteristically, such devices release the active agent at a high initial rate which then drops off in the manner of a first order equation.\nThe release of an active agent from a uniformly dispersed agent in a non-swellable homogeneous and imperforate plasticized polymer matrix by active agent permeation can be classified as a classical Fickian release, since the release pattern generally follows Fick's law of diffusion; EQU M.sub.t =Kt.sup.1/2\nwhere M.sub.t is the amount active agent released, K is a constant and t is time.\nU.S. Pat. No. 3,923,939 describes the removal of a portion of active agent from the surface of the aforementioned prior art diffusion and osmotic bursting devices by washing to form a depleted layer of polymer matrix. However, the devices do not include hydrogels, and moreover still generally exhibit a Fickian release, albiet with a lower initial release as evidenced by the release profile of the diffusion type devices of Example 1A with those of 1B in FIGS. 3 and 4 of that patent. The only anomolous release pattern described and exemplified is that of an osmotic bursting device in which the bursting device was washed for 24 hours such that the initial release rate was decreased by 90% as described in Example 3B and FIG. 10 of U.S. Pat. No. 3,923,939.\nOne drawback of the drug diffusion devices of the type exemplified in Example 1B of the aforementioned patent is that after treatment of the surface, the drug will diffuse at a rate corresponding to the drug diffusion rate in the device matrix until the drug is again uniformly distributed throughout the diffusion device. As a result, such a drug soluble rubbery polymer system exhibits numerous drawbacks.\nActive agent releasing devices containing a swellable hydrogel matrix are also known in the art. In such devices, the active agent is uniformly distributed in the non-swollen, or glassy, hydrogel polymer. Upon contact with an aqueous environment, the dry hydrogel swells as water penetrates the glassy matrix. The boundary between the glassy phase and the swollen, or rubbery, phase of the hydrogel is known as the solvent front. As the aqueous solvent front moves inward from the surface to the center of the hydrogel device, the active agent in the rubbery swollen phase of the hydrogel dissolves and diffuses through the swollen phase into the external aqueous environment. The active agent is substantially incapable of diffusion to any significant extent in the glassy non-swollen phase of the hydrogel. The rate of active agent transport into an aqueous environment is dependent upon a number of factors including the rate of penetration of the aqueous solvent front, the shape of the device, the diffusivity of the active agent in the swollen phase, the amount of active agent loading in the hydrogel matrix, the distance between the solvent front and the surface of thw device, the rate of decrease of the surface area of the glassy phase during solvent penetration and the like. See, for example, Ping I. Lee, Polymer Comunications, Vol. 24, p. 45-47 (1983).\nCharacteristically, the release rate for such hydrogel devices, containing uniformly distributed active agent, generally conforms to the following equation EQU M.sub.t =Kt.sup.a\nwhere M.sub.t is the amount of active agent released, K is a constant, t is time and a is between about 0.8 and 0.5. As a result, such systems do not provide for a substantially zero-order release, since the rate of release continuously drops off from a high initial release rate."}
{"text": "A wide variety of coatings have been used to coat the surfaces of packaging articles (e.g., food and beverage cans). For example, metal cans are sometimes coated using “coil coating” operations, i.e., a planar sheet of a suitable substrate (e.g., steel or aluminum metal) is coated with a suitable composition and cured. The coated substrate then is formed into the can end or body. Alternatively, liquid coating compositions may be applied (e.g., by spraying, dipping, rolling, etc.) to the substrate and then cured.\nPackaging coatings should preferably be capable of high-speed application to the substrate and provide the necessary properties when cured to perform in this demanding end use. For example, the coating should be safe for food contact, have excellent adhesion to the substrate, and resist degradation over long periods of time, even when exposed to harsh environments.\nMany current packaging coatings contain mobile or bound bisphenol A (“BPA”) or aromatic glycidyl ether compounds or PVC compounds. Although the balance of scientific evidence available to date indicates that the small trace amounts of these compounds that might be released from existing coatings does not pose any health risks to humans, these compounds are nevertheless perceived by some people as being potentially harmful to human health. Consequently, there is a strong desire to eliminate these compounds from food contact coatings.\nFrom the foregoing, it will be appreciated that what is needed in the art is a packaging container (e.g., a food or beverage can) that is coated with a composition that does not contain extractible quantities of such compounds. Such packages, compositions and methods for preparing the same are disclosed and claimed herein."}
{"text": "1. Field of the Invention\nThe present invention relates to a speed change control method for an automotive automatic transmission.\n1. Description of the Prior Art\nAn automotive automatic transmission is generally provided with a controller for determining, in accordance with a shift pattern set beforehand, a gearshift position suitable to a running condition of a vehicle, based on detected pieces of information such as throttle opening, vehicle speed, and a gearshift position currently established. The transmission is operable to automatically establish an optimum gearshift position under the control of the controller. In a typical shift pattern, a group of upshift lines for upshift operations from lower gearshift positions to higher gearshift positions, and a group of downshift lines for downshift operations from higher gearshift positions to lower gearshift positions, are set as a function of throttle opening and vehicle speed, as exemplary shown in FIGS. 1 and 2. When the vehicle running condition changes from a region indicated by numeral \"1\" in FIG. 1 to a region indicated by numeral \"2,\" for instance, an upshift operation from the first speed to the second speed takes place. Also, a downshift operation from the fifth speed to the fourth speed takes place when the vehicle running condition changes from a region indicated by numeral \"5\" to a region indicated by numeral \"4,\" for instance. Likewise, an upshift or downshift operation between other regions also takes place.\nIn this manner, the shift pattern is set such that a gearshift position for higher speed is selected at a lower vehicle speed region if the throttle opening is small, whereas a gearshift position for lower speed is selected in a higher vehicle speed region if the throttle opening is large. Thus, even if a vehicle runs on a descending road, a gearshift position for higher speed by comparison with vehicle speed is selected when the throttle opening is made small, so as to operate an engine in a lower rotational speed region. That is, an upshift can take place when the throttle is restored during the vehicle running on a descending road. Thus, engine braking is sometimes does not work effectively when the vehicle runs on a descending road. In this case, the vehicle speed increases, and hence a braking force is needed, so that a driver is required to operate a foot brake, or operate a select lever for switching drive ranges so as to cause engine braking to exercised.\nTo obviate such a drawback, an automatic transmission has been proposed, in which a downshift takes place when it is determined that a vehicle runs on a descending road. Such a determination is made by comparing gradient resistance to the vehicle derived from detected information with a determination reference value.\nAccording to the proposed automatic transmission, however, since the determination reference value for the determination of vehicle running on a descending road is fixed, a downshift takes place uniformly when such a descending running determination condition is satisfied. Thus, application of engine braking is always initiated upon fulfillment of the determination condition. On the other hand, a vehicle running condition in which engine braking for vehicle travel on a descending road should be exercised varies depending on drivers. Thus, it is difficult to satisfy a preference of an individual driver by means of such a uniform downshift control.\nAs another countermeasure for eliminating the need of a foot brake operation or a select lever operation during the vehicle travel on descending roads, a fuzzy speed change control based on vehicle speed, road gradient, etc., indicative of a vehicle travel condition, has been proposed. The fuzzy speed change control makes it possible to automatically effect a downshift so as to exercise engine braking when a vehicle runs on descending roads.\nHowever, according to a conventional fuzzy speed change control designed to conduct fuzzy inference for shift pattern selection directly from variables indicative of a vehicle running condition, such as gradient and vehicle speed, fuzzy rules must be set for all the combinations of the vehicle running condition variables, to make it possible to select an optimum shift pattern even if a vehicle runs in any condition. In addition, optimally setting a large number of fuzzy rules is very difficult. On the other hand, if fuzzy rules solely for typical running conditions are set, a shift pattern selection can be made inappropriately depending on the vehicle running condition.\nAs mentioned above, a shift pattern includes an upshift line group and a downshift line group which are set as a function of throttle opening and vehicle speed. In this respect, if a speed change control is conducted in accordance with a single shift pattern, then an upshift or a downshift takes place uniformly each time an operating point represented by throttle opening and vehicle speed crosses an upshift line or a downshift line. This sometimes makes it difficult to perform a speed change operation suitable to vehicle driving on ascending roads or a speed change control demanded by a driver.\nIn order to achieve a speed change control suitable to vehicle driving on ascending roads, a speed change control apparatus has been proposed in which a shift pattern for flat roads and a shift pattern for ascending roads which makes an upshift difficult to take place, are set beforehand. Further, a vehicle running on an ascending road is determined to select the shift pattern for ascending roads, if several variables, such as vehicle speed and gradient, indicative of the vehicle running condition satisfy a particular condition. This thereby effects a speed change operation suitable to the vehicle running on ascending roads, while preventing an unnecessary upshift.\nHowever, in the proposed speed change control apparatus, the same shift pattern for ascending roads is used for both of a steep ascending road and a gentle ascending road. Depending on the gradient or the vehicle speed, this can prevent an upshift even if a driving force is sufficient, and can permit an upshift even if a driving force is deficient, resulting in shift hunting.\nFurther, to achieve a speed change operation demanded by a driver, a speed change control apparatus has been proposed in Japanese Provisional Patent Publication No. 2-275174, in which an ordinary shift pattern (normal pattern or economy pattern), and a sporty pattern (power pattern) where shift lines are provided on the higher speed side than those of the normal pattern are set beforehand, and in which an arbitrary one of these patterns is selectable manually or automatically. According to a speed change control apparatus of this kind, changeover from the normal pattern to the sporty pattern automatically takes place when an accelerator depression speed exceeds a threshold value, which varies dependent upon the vehicle speed and the throttle opening, under a condition where predetermined requirements such that a select lever is in a D-range are fulfilled.\nThe threshold value of the accelerator depression speed is set, as exemplary shown in FIG. 3. In the case of this setting example, the entire vehicle running area is divided into sixteen regions by use of throttle opening and vehicle speed as parameters, and threshold values T11 through T44 for the respective regions are set. The threshold values T11 through T44 are set such that one for a higher vehicle speed region has a larger value than one for a lower vehicle speed region (T11&lt;T12&lt;T13&lt;T14, for instance), so that the switching to the sporty pattern is likely to take place at a lower vehicle speed. Further, the threshold values T11 through T44 are set such that one for a larger throttle opening region has a smaller value than one for a smaller throttle opening region (T11&gt;T21&gt;T31&gt;T41, for instance), so that the switching to the sporty pattern is likely to take place at a larger throttle opening. As a consequence, a shift from the normal pattern shown by solid lines in FIGS. 1 and 2 to the sporty pattern shown by doted lines takes place when the accelerator pedal is depressed at a speed higher than an associated threshold value (T11, - - - , or T44).\nHowever, the proposed apparatus, which only permits selection of either one of the two shift patterns, cannot achieve an elaborated speed change operation which fully meets a vehicle running condition or a driver's demand.\nAs another prior art for shift pattern change, it has been also known to set a plurality of, e.g., five upshift characteristic lines SC1 through SC5 in connection with an upshift from the third speed to the fourth speed, for instance, as exemplary shown in FIG. 4, and select one of these shift characteristics in accordance with a vehicle running condition represented by throttle opening, engine rotational speed, vehicle speed, lateral acceleration, and longitudinal acceleration, for instance.\nHowever, using a large number of shift patterns in this manner requires a lot of labor in setting the shift patterns, and takes a lot of time for determination processing to select an optimum pattern from the shift patterns which is large in number, causing a delay in shift control to a change in the vehicle running condition.\nFurther, to eliminate drawbacks produced when a speed change control is effected with use of a conventional shift schedule which requires manual speed change from the D-range to the 2nd-range or a brake pedal depressing operation to compensate a deficiency in engine braking performance which occurs when a vehicle runs on descending roads, it has been proposed in Japanese Patent Publication No. 52-20630 to use a particular shift schedule including a specific speed region, corresponding to a nearly fully closed throttle opening, for automatic speed change, the specific speed region being expanded toward a predetermined vehicle speed which is usually included in an adjacent speed region for which a gearshift position which is one level higher than that for the specific speed region is established (see, FIG. 47).\nHowever, In case that a downshift line for a small throttle opening region is set on the higher speed side, as in this proposal, even if a driver wishes to make a coasting vehicle running, a downshift takes place contrary to the driver's intention, to cause engine braking to be exercised. This causes the driver to have a feeling of dissatisfaction or incompatibility.\nOn the other hand, in case that a downshift line for a small throttle opening region is set on the lower speed side, as distinct from the aforementioned proposed shift schedule, a downshift does not take place as long as the accelerator pedal is not depressed by the driver. Thus, a delay can be caused in downshift, e.g., when the vehicle passes a sharp corner in sporty driving during which the driver makes braking at a location short of the corner and then depresses the accelerator pedal for re-acceleration. That is, the downshift does not take place at the location short of the corner, but takes place when the vehicle reaches the exit of the corner. In this case, acceleration at the exit of the corner is prevented by the downshift effected at the corner exit.\nAs explained in the above, it is difficult heretofore to achieve an elaborated speed change control with the use of a relatively simplified control system."}
{"text": "This invention relates to portable electronic video recording devices.\nIn the past, video cameras have been bulky, and those in most widespread use have been hand-held or, in some cases, shoulder mounted. Such devices have been prone to producing shaky pictures as a result of the unsteadiness of the operator's hands. Also, it has been difficult to film a subject without that subject's being aware that he is being recorded. As a result, the subject tends to behave other than normally, due to self-consciousness."}
{"text": "Many vehicles are used over a wide range of vehicle speeds, including both forward and reverse movement. Some types of engines, however, are capable of operating efficiently only within a narrow range of speeds. Consequently, transmissions capable of efficiently transmitting power at a variety of speed ratios are frequently employed. When the vehicle is at low speed, the transmission is usually operated at a high speed ratio such that it multiplies the engine torque for improved acceleration. At high vehicle speed, operating the transmission at a low speed ratio permits an engine speed associated with quiet, fuel efficient cruising. Typically, a transmission has a housing mounted to the vehicle structure, an input shaft driven by an engine crankshaft, and an output shaft driving the vehicle wheels, often via a differential assembly which permits the left and right wheel to rotate at slightly different speeds as the vehicle turns.\nIn some types of transmissions, a gear ratio between the input shaft and the output shaft is established by engaging one or more friction clutches. During some operating conditions, a clutch may be partially engaged such that it transmits torque between two elements that are rotating at different speeds or between a rotating element and a stationary element such as the transmission housing. For example, during an upshift, the oncoming clutch may be partially engaged while the engine speed gradually declines to the speed associated with the final gear, enabling the transmission to provide output torque during the shift event. As another example, a clutch may be partially engaged while the vehicle accelerates from stationary to the speed associated with complete engagement of first gear. When a clutch is partially engaged, power is converted into heat which must be absorbed by the clutch and eventually dissipated.\nA common type of clutch utilizes a clutch pack having separator plates splined to a rotating housing and interleaved with friction plates splined to a rotating shell. When the separator plates and the friction plates are forced together, torque may be transmitted between the housing and the shell. Typically, a separator plate on one end of the clutch pack, called a reaction plate, is axially held to the housing. A piston applies axial force to a separator plate on the opposite end of the clutch pack, called a pressure plate, compressing the clutch pack. The piston force is generated by supplying pressurized fluid to a chamber between the housing and the piston. Fluid typically flows between the separator plates and friction plates keeping the friction material on each friction plate damp. Heat is generated at the interface when the clutch is partially engaged. The heat is absorbed by the separator plates and transferred to the fluid by convection and is removed from the clutch with the fluid."}
{"text": "Typically, real user tracking and monitoring of web applications is accomplished by embedding monitoring agents into web pages. During navigation of a web page, the monitoring agent captures performance metrics related to the web page and sends data back to a server for analysis and reporting. On the other hand, synthetic testing tools rely on customization of a user's web browser to capture performance metrics related to web pages. Because the collection mechanisms differ between real user monitoring and synthetic testing, the data collected differs and cannot be easily compared.\nTherefore, it is desirable to develop a technique for synthetic testing of web applications which results in performance metrics that are the same as the performance metrics collected in real user monitoring and collected in the same manner as the performance metrics collected in real user monitoring.\nThis section provides background information related to the present disclosure which is not necessarily prior art."}
{"text": "Unless specifically indicated otherwise, the materials described in this section are not prior art to the claims in this application, and are not admitted to be prior art by inclusion in this section.\nThe material in this patent document is subject to copyright protection under the copyright laws of the United States and other countries. The copyright owner grants a limited license to any member of the public to reproduce the patent document as it appears in official governmental records. All other copyrights rights are reserved.\nMarine anchors are commonly used to secure a marine vessel to a mooring bed. A vessel may utilize one or more anchors for securing the bow of the vessel, and one or more anchors for securing the stern of the vessel. When a marine anchor is not in use, the anchor must be stowed on the vessel.\nSome vessels may be equipped with anchor lockers or lazarettes that are designed specifically for the purpose of stowing an anchor, or that may be adapted to that purpose. Anchor lockers and lazarettes are compartments built beneath the hull of a vessel, and are generally accessible by means of a door or a hatch. Anchor lockers and lazarettes have the disadvantage that they do not provide for convenient access to the anchor. The door or hatch must be opened in order to access the stowed anchor, and the user must reach down into the anchor locker or lazarette to access the anchor and its rode. Another disadvantage of an anchor locker or lazarette is difficulty of use. The anchor and rode must be lifted out from its stowage location below the deck and carried to the position on the vessel where the anchor is to be deployed, typically the bow or the stern of the vessel. A further disadvantage of anchor lockers and lazarettes, which are built into the structure of the vessel, is that the anchor locker or lazarette may have a size limitation or shape factor that restricts the size or style of anchor that may be accommodated by the fixed, built in anchor locker or lazarette.\nAnchor stowage devices have been developed as alternatives to anchor lockers and lazarettes, for vessels that do not have anchor lockers or lazarettes and/or to overcome the disadvantages of anchor lockers and lazarettes.\nOne category of conventional anchor stowage devices, commonly called anchor davits, are designed to project beyond the deck and hull of a vessel and to stow the anchor in a horizontal orientation. These anchor davit devices are typically intended for use at the bow of a vessel. A representative anchor davit design is disclosed in U.S. Pat. No. 5,996,524 to Johnson. Johnson discloses an anchor davit comprising an elongated U-shaped frame for receiving a main shank of an anchor; a roller at the outboard end of the frame for guiding the anchor line as the anchor is lowered or raised; a smooth round rod means at the outboard end of the frame for providing chafe protection of an anchor line; and means at a rear portion of the U-shaped frame for securing the frame to a bow, a deck, a gunnel, or another structural member of a boat. Anchor davit devices share the disadvantage that they take up deck space to the extent that the devices are mounted on the deck, and/or they increase the effective length of the vessel to the extent that they extend forward of the bow of the vessel. This disadvantage renders davit-style anchor stowage devices particularly undesirable for stowage of anchors at the stern of the boat, where deck space in the cockpit area is at a premium and where protrusions outboard of the hull and/or deck are inconvenient. A further disadvantage of davit-style anchor devices is the absence of a means for securely holding the shaft of an anchor in a manner that prevents shifting and rattling of the anchor.\nOther anchor stowage designs attempt to address the shortcomings of the davit-style anchor devices, especially for use with stern anchors. However, each of these designs possess disadvantages of their own.\nU.S. Pat. No. 4,479,452 to Awalt discloses a storage device for a swivel-fluke-type marine anchor comprising a shank-receiving boom and a fluke channeling member both being rotatably mounted to the stern of a vessel, in combination with a back plate fixedly mounted to the stern of the vessel. To capture a swivel-fluke-type marine anchor, the shank-receiving book and fluke channeling member are rotated to extend horizontally outboard from the stern of a vessel. Once captured, the anchor is stowed by rotating the shank-receiving book and fluke channeling member to a vertical stowage position against the back plate. One obvious disadvantage of the invention of Awalt is that it is designed to accommodate one particular anchor design type: a swivel-fluke-type marine anchor. This disadvantage, the lack of compatibility with a variety of marine anchor types, is typical of most conventional marine anchor stowage devices. Another disadvantage of the invention of Awalt is complexity. This design complexity drives complexity of use. To release an anchor, rode is released from a bit arm, a locking pin must be removed from a locking arm. A control force lever arm must be depressed causing the fluke receiving channel to project outwardly away from the vessel, and the frame rotates to its deployed position. The anchor is then in position to be deployed. A similar level of complexity is involved in capturing and stowing the anchor. In addition to driving complexity of installation and use, design complexity drives cost.\nU.S. Pat. No. 4,164,194 to Kurz discloses an anchor stowage device that is designed to suspend an anchor above the deck of a vessel. The invention of Kurz comprises a base support affixed to the deck of the vessel and a bracket attached to a railing above the deck. To stow a swivel-fluke-type marine anchor, the stabilizing bar of a swivel-fluke-type anchor is inserted into channels in the base support, the shank of the anchor is then inserted into a channel in the bracket, a locking feature is rotated to close the opening of the channel in the bracket, and a wing nut is tightened onto a threaded stud in order to secure the locking feature such that the anchor cannot come out of the channel. The invention of Kurz is limited to use with a swivel-fluke-type marine anchor. A further disadvantage of the invention of Kurz is that the stowed anchor takes up deck space. An additional advantage is the need to release the wing nut in order to deploy an anchor, and to install and tighten the wing nut in order to stow an anchor. This unfastening/fastening step is inconvenient, particularly where quick deployment of an anchor is desired, and the lose wing nut hardware is subject to being dropped into the vessel or overboard.\nAs is readily apparent from a review of conventional anchor stowage device designs, what is needed is an anchor stowage device that is simple in design and use and may be rapidly deployed, that minimizes use of valuable deck space, and is capable of accommodating a variety of marine anchor types, including, but not in a limiting sense, fluke-style anchors, claw-style anchors and plow-style anchors."}
{"text": "Lower leg exercise involves muscles from the knee to the ankle and muscles surrounding the ankle. Lower leg muscles are exercised by rotational movements of the foot about the ankle joint. Lower leg exercise machines commonly provide one or two footplates with adjustable resistances against which the user may exert a force. Such exertion, when regularly repeated, results in lower leg muscle strengthening, improved muscle endurance, and increased joint motion flexibility. These results are key to prevention and rehabilitation of lower leg muscle and joint injuries.\nAs physical exercise for recreation and for cardiovascular fitness increases in popularity, so do the injuries associated with the lower leg. Ankle sprains are the most frequent injuries associated with recreation and fitness exercise. Knee injuries are also common. Running and jumping, for example, are known to cause a knee injury medically described as infrapatellar tendinitis, resulting in pain at the shin just below the kneecap. This particular injury and its rehabilitation exercises are described in an article \"How I Manage Infrapatellar Tendinitis\" by Joseph E. Black, MD, in the October, 1984, issue of The Physician and Sportsmedicine, pages 86-92. A principal object of the present invention is to enable individuals to perform the recommended exercise to prevent and rehabilitate this injury.\nWhen serious athletes have lower leg injuries, the supervised care of an athletic trainer or physical therapist may be appropriate to hasten the return of the athlete to full competitive ability. However, for the non-athlete, fitness buff, or recreational player, the objective is not the fastest return to competition. It is instead to quickly regain the capability of walking straight ahead on both feet without a limp. Such movement involves the feet pivoting back and forth in a substantially vertical plane. This motion is technically known as dorsiflexion and plantarflexion.\nCommercially available lower leg exercise machines are designed for use at centralized facilities, such as sports medicine clinics and university training rooms. These machines provide for foot movement against resistance in two or three different planes simultaneously. Some have built in seats and elaborate resistance measurement feedback systems. They are complex and costly machines that are economically justified by application to multiple patients. Examples include U.S. Pat. No. 4,650,183 to McIntyre, which discloses a two axis machine with seat and hydraulic gage console; U.S. Pat. No. 4,733,859 to Kock et al, which shows a two axis machine with seat; and U.S. Pat. No. 4,452,447 to Lepley et al., disclosing a three axis machine. These machines provide for exercise in the inversion/eversion and adduction/abduction planes in addition to the dorsiflexion/plantarflexion plane in order to provide complete ankle rehabilitation. Missing from the marketplace are affordable, single-axis, lower leg exercisers for individuals to use at their own convenience at home.\nIn hospitals where patients are bedridden, a nurse often must manually massage a patient's lower legs to stimulate circulation so that varicose veins and the development of thrombophlebitis can be avoided. There is a need for a small and light weight lower leg exerciser that can be attached by a nurse to a supine patient's foot, with means for providing exerciser stability on a bed, so the patient can perform a lower leg exercise by him or herself. Such an exerciser might eliminate the need for manual massage since exercise stimulates blood circulation. Also, leg exercise would help to maintain lower leg muscle strength in the immobile patient.\nU.S. Pat. No. 4,159,111 to Lowth, U.S. Pat. No. 3,525,522 to Piller, and U.S. Design Pat. No. 189,011 to Berne disclose single-axis lower leg exercisers. These machines have significant deficiencies for supine patient use, however. Because the footplate pivot points are located opposite the footplate from the user's ankle joint, substantial movement of the patient's leg is required to rotate a footplate. This movement is easily accomplished in a sitting position with the knee bent, but it is difficult in the supine position with the leg straightened. Also, these inventions are shown with horizontal base plates beneath the footplate pivots. The bases of the Lowth, Piller and Berne exercisers would be unstable to torsional moments when resting on a non-rigid bed surface. They would slide and tilt on the bed surface.\nResistance exercise of major body muscle groups is commonly done using free weights grasped by the hands or by using machines which have cables connecting pivoting members to weight stacks. These resistance means are not practical for ankle exercise. The high resistance levels needed for major muscle groups are not needed for ankle exercise. And the relatively small angle of rotation of the foot (short stroke) permits more compact resistance means to be used. Hydraulic cylinders and friction disks are the predominant resistance means used for lower leg exercisers. U.S. Pat. No. 4,605,220 to Troxel shows four hydraulic shock absorbers attached between a footplate and base. U.S. Pat. No. 4,650,183 to McIntyre shows two rotary hydraulic actuators connecting a footplate to a base. And U.S. Pat. No. 4,452,447 to Lepley et al. shows three hydraulic cylinders connected from a base to footplate cranks. Although smooth and speed controlled motion is provided by hydraulic resistance, disadvantages are its high cost and eventual leakage of hydraulic fluid. Friction disks, on the other hand, are inexpensive, more compact, and don't leak. The latter two features are especially important for a portable device for use in a hospital bed, for example.\nMost exercise machines provide a constant resistance throughout the range of motion of a body member about its joint. But as the body member moves, the leverage of the muscle(s) being stressed changes. There is a position at which leverage is greatest. At the ends of the range of body member motion the leverage is usually lowest. Thus, if the resistance is set for optimum muscle stress at the greatest leverage position, the range of motion under this resistance will be very limited. On the other hand, if the resistance is set for optimum muscle stress at the ends of the range of motion, the resistance may be too low at the center of the range of motion. It is desired to have the resistance level vary with muscle leverage so that each muscle is evenly stressed throughout its range of movement."}
{"text": "This invention relates to methyl-thio-2-methyl-2-pentenoate derivatives which are defined according to one of the structures: ##STR3## (representative of cis and/or trans isomers) wherein the wavy lines represent moieties juxtaposed around the carbon-carbon double bond, and uses thereof in augmenting or enhancing the aroma or taste of foodstuffs, chewing gums, medicinal products or toothpastes.\nArtificial flavoring agents for foodstuffs have received increasing attention in recent years. In many areas, such food flavoring agents have been preferred over natural flavoring agents at least in part due to their diminished cost and their reproducible flavor qualities. For example, natural food flavoring agents such as extracts, concentrates and the like are often subject to wide variations due to changes in the quality, type and treatment of the raw material. Such variations can be reflected in the end product and result in unfavorable flavor characteristics in said end product. Additionally, the presence of the natural product in the ultimate food may be undesirable because of increased tendency to spoil. This is particularly troublesome in food and food uses where such products as dips, soups, chips, sausages, gravies and the like are apt to be stored prior to use.\nThe fundamental problem in creating artificial flavor agents is that the artificial flavor to be achieved be as natural as possible. This generally proves to be a difficult task since the mechanism for flavor development in many foods is not completely known. This is noticeable in products having cocoa, chocolate, liver, ripe fruit and strawberry flavor characteristics.\nReproduction of cheesy, surface ripened cheese-like, cocoa powder-like, liver-like and fruity flavor and aroma nuances has been the subject of long and continuing searches by those engaged in the production of foodstuffs and beverages. Thus, severe shortages of food in many parts of the world has given rise to the development of previously unused sources of protein which are unpalatable. Accordingly, the need has arisen for the use of flavoring materials having such aroma and taste nuances which will make such sources of protein palatable to human sensory organs.\nThioesters in particular are well known for use in flavoring of foodstuffs. Thus, U.S. Pat. No. Re. 30,370 discloses compounds having the generic structure: ##STR4## wherein R.sub.1 and R.sub.2 are alkyl or alkenyl, for example, allyl thiopropionate having the structure: ##STR5## are disclosed as being useful for meaty flavors.\nUnsaturated alkenoic acid esters (containing the oxygen rather than the sulfur atom) are also known in augmenting or enhancing the aroma or taste of strawberry flavored foodstuffs, for example, those disclosed in U.S. Pat. No. 3,931,293, for example, the compound having the structure: ##STR6## and corresponding trans isomers and mixtures of same having a high \"cis\" isomer content. Nothing in the prior art discloses the novel compounds of our invention, the methyl-thio-2-methyl-2-pentenoates of our invention or the unexpected, unobvious and advantageous organoleptic uses of same."}
{"text": "The present invention concerns that of a new and improved swing covering accessory. The accessory is used to cover an outdoor infant/toddler swing and is designed to be placed into the swing and around the outside of it, providing both entertainment and sanitary purposes for infants and toddlers using the swing."}
{"text": "Conventionally, as is shown in FIGS. 6 or 7, an automobile seat is constructed by a seat cushion 1, seatback cushion 2 and seat adjuster unit 3. When assembling, seat cushion 1 and seatback cushion 2 are independently formed in order to cover the cushions 1 and 2 by separate cover materials. As shown in FIG. 7, cushion material 5 of the seat cushion 1 is fixed on a seat cushion pan 4 and cover material 6 is put on the cushion material 5. The end portion 7 of the cover material 6 is put under the cushion pan 4 to be secured thereto by clip means 8.\nAs shown in FIG. 6, seatback cushion 2 includes cushion material and seatback frame (not shown) secured to the cushion material. Seat adjuster 3 is secured to the cushion and cover material 9 is attached over the seatback frame by clip means at an end portion 10 of the cover material 9. Cover 11 is finally attached to the assembly to form the seatback cushion 2. Seat cushion 1 is secured to the seatback cushion 2 by bolt means 12 through seat adjuster 3 which is integrally secured to the seatback cushion 2.\nThis conventional seat assembly, however, has a drawback that since the connection between the two cushion portions are made through bolt and pin connection at the relatively thin seat cushion pan, load concentration at the connecting portion will cause to reduction of rigidity and strength. If a reinforcement is attached to the connecting portion in order to cover the strength or rigidity, weight will be increased, which leads to the cost increase."}
{"text": "A sagittal saw is a powered surgical tool that is often used in an orthopedic surgical procedure. A sagittal saw generally includes a handpiece that houses a motor and the complementary control circuit that regulates the actuation of the motor. Extending forward, distally, from the handpiece, is a planar saw blade. The most forward end of the saw blade is formed with teeth for cutting hard tissue against which the blade is applied. A drive mechanism internal to the housing transfers the power developed by the motor to the blade. More particularly, the drive mechanism converts the rotary motion produced by the output shaft of the motor to the blade so that the blade moves in an oscillatory, back-and-forth pattern in the plane in which the blade is aligned. Consequently, when a sagittal saw is actuated, the blade teeth move in a back-and-forth pattern against the hard tissue or bone to which the teeth are applied. As a consequence of this motion and the forward pressure applied by the surgeon holding the saw, the teeth cut and separate the hard tissue or bone.\nA sagittal saw is often used in an orthopedic surgical procedure to selectively remove bone. One particular type of orthopedic surgical procedure in which the saw is used is a joint replacement procedure. As the name implies, in this type of procedure, the surgeon resects the bone between the joints in the patient and substitutes an artificial joint.\nIn orthopedic surgical procedures, it is very important to ensure that, when a bone section is separated from the rest of the bone, the section is removed along very precise cut lines. This is very important in a joint replacement procedure because the substitute joint typically has a component that is designed to precisely fit in the space defined by cut lines of the section of bone that is left in place.\nTherefore, in order to insure that the proper cut lines are formed in the bone, the surgeon typically first mounts a cutting guide, sometimes called a jig, to the bone adjacent to the location where the cut is to be made. One type of cutting guide is in the form of a block formed with a precisely shaped set of slots. The slots define the lines along which the bone is to be cut. The surgeon then performs the surgical procedure by sequentially inserting the saw blade in the slots. Once the blade is inserted in a slot, the saw is actuated. In this manner the surgeon is able to cut the bone along the precisely defined lines along which the bone is to be separated.\nAnother type of cutting guide is in the form of an open-face block which is mounted to the bone at the appropriate location and defines a guide surface thereon. The surgeon performs the surgical procedure by positioning the saw blade flat against the guide surface to make the desired cut. This type of cutting guide is advantageous in that the visibility of the cutting area is improved from that of cutting guides utilizing slots, since the saw blade is not hidden within the slot.\nWhile presently available sagittal saws and complementary cutting guides work reasonably well, there are some noticeable limitations. As mentioned above, known commercially available sagittal saws are provided with flat planar blades that oscillate. This type of blade invariably rubs against the cutting guide material forming the slot(s) in which the blade is inserted. This repetitive contact wears away this slot-defining material. One problem associated with the wearing away of this material is that it widens the slot. Eventually, the slot may become so wide that it no longer serves to precisely define the cut line it is intended to define. Once a cutting guide is so worn, it needs to be replaced. Moreover, the wearing of the material forming the cutting guide generates a fine dust of the material. This dust inevitably settles on the surgical site at which the procedure is being performed. Consequently, during the procedure, the surgical personnel are required to spend an appreciable amount of time flushing the site to remove this dust. Having to repeatedly perform this process runs counter to one of the primary goals when performing surgery, that one should perform the procedure as quickly as possible to both minimize the likelihood the exposed tissue is open to infection and the amount of time the patient is held under anesthesia.\nAs discussed above, the oscillating blade of the current surgical saw will repeatedly gall the surfaces of the cutting guide forming the slot in which the blade is inserted. One further disadvantage of this blade galling is that same consumes power, and since many sagittal saws are battery powered, the power expended due to the friction caused by the blade galling reduces the overall amount of power the battery has available to power the saw. In other words, the power consumed to overcome this frictional contact can reduce the overall amount of time the battery, on a single charge, is able to power the saw. Moreover, as a consequence of the saw blade galling against a surface of the cutting guide, and then pulling away from the surface, there is some jerking of the blade. This jerking motion is transferred from the blade through the handpiece into the hand of the surgeon holding the saw. Consequently, the surgeon must exert some muscle control to hold the handpiece steady when he/she is exposed to this jerking motion.\nAlso, as an inevitable result of the back-and-forth motion of the blade, the surgical saw invariably vibrates. Again, the surgeon is required to engage in some conscious or unconscious physical effort to hold the saw steady when it so vibrates. Over time, having to so hold the saw steady to overcome this vibration can be significantly mentally and physically fatiguing.\nRecently, there has been proposed a new type of surgical sagittal saw that does not include a flat oscillating blade. This saw instead includes an endless metal band that is formed with outwardly directed teeth. The band is wrapped around a static guide that extends forward from the handpiece, with the band teeth extending outwardly from the guide. A drive mechanism rotates the band. Since the guide of the proposed saw does not move, it is believed that many of the problems associated with saws that have oscillating blades would be eliminated. One example of this type of saw is disclosed in U.S. Pat. No. 5,725,530.\nHowever, there are disadvantages associated with the above saw. It is expensive to provide the toothed metal band. Also, the metal band appears prone to fatigue and, consequently, breakage. The time it takes to replace this metal band while in the middle of a surgical procedure can appreciably increase the overall time required to perform the procedure.\nMoreover, one of the goals of modern surgery is to, whenever possible, perform the procedure using a minimally invasive surgical (MIS) practice. As the name implies, in an MIS procedure only a relatively small incision is made with minimal soft tissue disruption in order to gain access to the surgical site. Minimizing the extent to which a patient's tissue is exposed reduces the amount of tissue that is exposed to the ambient environment and the potential for infection caused by such exposure. Furthermore, reducing the extent to which the patient's tissue is incised minimizes the amount of tissue that then needs to heal.\nIn order to perform a minimally invasive surgical procedure on a bone or bone joint, only a relatively small portion of the surrounding soft tissue is incised to expose the bone or the joint. Consequently, the bone or joint is not well exposed. In order to make the resection in the bone, the oscillating saw blade that is employed is typically longer than the saw blade used to perform a conventional, resection surgical procedure. Given the relatively long length of the blade, the blade has a mass moment of inertia that is appreciably greater than the mass moment of inertia associated with shorter-length blades. Consequently, when the saw to which this blade is attached is actuated, more vibratory motion is created by the blade and transferred to the rest of the handpiece than when a shorter blade is used. This increased vibratory motion can make it difficult for the surgeon to hold the saw steady. Further, because this longer type of saw blade oscillates along essentially its entire length during the cutting procedure, the saw blade can cause extensive damage to the soft tissue at the incision.\nMoreover, a longer blade is more flexible than its shorter counterpart. The added flexibility of the blade can result in the blade making less precise cuts in the bone the blade is used to shape. Unfortunately, it is not possible to reduce this flexibility by simply increasing the overall thickness of the blade. Taking this action increases the mass of the blade and, by extension, the mass moment of inertia of the blade. For the reasons set forth above, increasing the mass moment of inertia of the blade would, in turn, increase the extent to which the associated saw, when actuated, vibrates.\nAlso, both the sagittal saws that are provided with oscillating blades and the proposed saw with a static guide bar are used with cutting guides that are relatively large in size. The relatively large size of the presently available cutting guides makes it difficult, if not impossible, to perform a minimally invasive surgical procedure.\nU.S. Pat. No. 2,854,981 discloses a surgical saw having a saw blade pivotably supported at the end of a beam which extends forwardly from a handpiece. The beam includes a pair of tubes which are secured on opposite sides of a support rod, which tubes house therein reciprocating thrust rods. The thrust rods push against a plate adjacent the blade and cause the blade to undergo pivotal oscillating movement.\nThis saw, due to the pivotal movement of the blade at the distal end of the saw, would not appear to create excessive vibratory motion and/or soft tissue damage. However, the extension of the saw relative to the cutting slot or kerf created by the blade in the hard tissue or bone is limited. That is, the depth at which the saw can cut is limited to the length of the blade itself, since the beam which supports the blade is much larger than the blade.\nFurther, the arrangement illustrated in the above patent actuates the blade through compression of the respective thrust rods which then push against the blade to move same. This type of arrangement necessarily requires that the rods be dimensionally large and constructed of a heavy material capable of repeatedly withstanding these types of forces, which then results in a heavy and cumbersome saw."}
{"text": "This invention relates to a transistor circuit comprising an emitter-coupled transistor pair comprising a first and a second transistor, and a third transistor to compensate for the base-current of the first transistor, which third transistor has its collector-emitter path arranged in series with the collector-emitter path of one of the transistors of the transistor pair.\nSuch a transistor circuit is disclosed in U.S. Pat. No. 3,714,600. This prior-art transistor circuit comprises an emitter-coupled transistor pair, the base-current of the first transistor being compesated for by the base-current of a third transistor having its collector-emitter path arranged in the collector line of the first transistor. The base-current compensation causes the impedance at the base of the compensated first transistor to increase so that the load on the signal source to which said base in coupled is reduced substantially. The base-current of the second transistor can be compensated for identically.\nIn gyrator-capacitor filters, which are described , inter alia, in \"Integration of Analog Filters in a Bipolar Process\", IEEE Journal of Solid-State Circuits, Vol. SC-17, No. 4, Aug. 1982, pp. 713-722, the emitter-coupled transistor pair is an important element, which is often referred to as a transconductance circuit or transconductor, the voltage difference between the bases of the transistors of one pair being converted into two balanced output currents. By coupling two emittercoupled transistor pairs a gyrator can be formed. The quality and precision of the gyrator-capacitor filters depends inter alia on the manner in which a gyrator loads the other filter components, such as capacitors, resistors and other gyrators. In the ideal case a gyrator comprises two transconductors each having an infinitely high input impedance such that the transconductors neither load one another nor the filter components coupled to the gyrator. A high input impedance might be obtained by base-current compensation as used in the prior art transistor circuit. However, the prior-art transistor circuit requires a current mirror for the purpose of base-current compensation. The transistors of this current mirror are of a conductivity type opposite to that of the transistor in the imetter-coupled transistor pair. In integrated circuits such transistor pairs are formed by NPN transistors on account of thier superior signal-amplifying properties. This means that PNP transistors will be needed for the current mirror. The currents in the PNP current mirror should be most faithful replicas of the base current of the transistor to be compensated. However, PNP transistors have a lower current gain and an inferior high-frequency performance in comparison with NPN transistors. Therefore, a current mirror comprising PNP transistors is less accurate and less suitable for use at high frequencies."}
{"text": "1. Field of the Invention\nThe present invention relates to a level converting circuit, and particularly to a signal-level converting circuit for bidirectionally transmitted logic signals.\n2. Description of the Related Art\nWhen a signal is transferred between circuits having different signal levels, a signal level converting circuit must be provided between the circuits. Conventionally, when only signal level conversion (hereinafter referred to as \"level conversion\") is intended between circuits, level conversion has been accomplished by an IC for level conversion provided between the two circuits. Such an IC, however, is capable of carrying out level conversion in only one direction of signal transfer between the circuits of differing signal levels. Consequently, in order to carry out level conversion bidirectionally between circuits such as bus lines which handle bidirectional signals, two types of level conversion ICs are necessary. Furthermore, switching between the bidirectional signals necessitates a number of switches corresponding to the signal lines making up the bus line as well as a control circuit for controlling switching of these switches. Hereinafter, level conversion in one direction of signal transfer between circuits of differing signal levels is referred to as one directional conversion, and level conversion in both directions of signal transfer between circuits of differing signal levels is referred to as bidirectional conversion.\nNow, an example of the above-described conventional bidirectional conversion circuit will be explained with reference to FIG. 1. FIG. 1 shows a block diagram of an example of a conventional bidirectional conversion circuit which bidirectionally converts a logic signal having a signal level of 5 volts (hereinafter referred to as a \"5-V logic signal\") and a logic signal having a signal level of 3 volts (hereinafter referred to as a \"3-V logic signal\").\nThe bidirectional conversion circuit is equipped with a first switch 32 to which each of 5-V input/output terminals ({1}-{8})31 is connected, a 5-to-3 level converter 33 that converts signal level from 5 volts to 3 volts, a second switch 34, 3-V input/output terminals ({1}-{8}) 35 connected to the second switch 34, a control circuit 36 for controlling switching of the first and second switches 32, 34, a control terminal 37 connected to the control circuit 36, and a 3-to-5 level converter 38 that converts signal level from 3 volts to 5 volts. The 3-volt power source terminal 40 and 5-V power source terminal 41 are connected to 5-to-3 level converter 33 and 3-to-5 level converter 38, respectively.\nFirst will be explained the operation of level conversion from 5 volts to 3 volts (hereinafter referred to as \"5-to-3 conversion\"). A 5-V logic signal is applied at 5-V input/output terminals 31, and after passing through the first switch 32, the 5-V logic signal is converted to a 3-V logic signal at 5-to-3 level converter 33. This logic signal is delivered to 3-V input/output terminals 35 by way of the second switch 34.\nNext will be explained the operation of level-conversion from 3 volts to 5 volts (hereinafter referred to as \"3-to-5 conversion\").\nControl circuit 36 is employed for 3-to-5 conversion. In this case, a synchronous clock signal, control data (information describing which switches to turn on and which switches to turn off), and a strobe signal that controls switching timing are supplied to control terminals 37. In response to the strobe signal, control circuit 36 outputs switching control signals corresponding to the control data. By this switching control signal, each switch of the first and second switches 32, 34 is switched, in this example, in the direction opposite that shown in the figure. The 3-V logic signals are received at 3-V input/output terminals 35, and after passing through the second switch 34, are converted to 5-V logic signals at 3-to-5 level converter 38. Subsequently, the level-converted signals pass through the first switch 32 and are delivered to the 5-V input/output terminals 31.\nRegarding the prior art signal level conversion, a circuit is disclosed in Japanese Patent Laid-open 227317 (92). This circuit is an output buffer circuit that converts a CMOS logic level to an ECL logic level. In this circuit, ECL voltage signals of high and low voltage levels are generated in advance by a reference voltage circuit. Signal level conversion from a CMOS level to a ECL level is effected by a switch that operates to connect either of said ECL voltage signals to an ECL load in accordance with the CMOS logic states.\nBecause the above-described bidirectional conversion circuit of the prior art can execute level conversion using an IC for one directional conversion, level conversion between circuits such as bus lines that handle bidirectional signals necessitates at least two types of level converters as well as switches in a number equal to the number of signal lines making up the bus lines and a control circuit unit for controlling the switches. Consequently, this type of bidirectional conversion circuit of the prior art has the drawback that a large number of parts and connectors are needed, resulting in complex circuit construction."}
{"text": "The conventional materials used as luminescent substrate include phosphor, nanocrystal, glass, etc. Comparing to the crystal and phosphor, the glass is transparent, rigid, and has excellent chemical stability and superior luminescent performance. In addition, the glass can be easily machined into products with various shapes, such as display devices or luminescent light sources with various shapes and sizes.\nFor example, in vacuum microelectronics, field emission devices usually use luminescent glass as illuminant, which has shown a wide prospect in illumination and display techniques and draws a lot attention to domestic and foreign research institutes. The working principle of the field emission device is that, in vacuum, the anode applies a positive voltage to the field emissive arrays (FEAs) to form an accelerating electric field, electron emitted from the cathode accelerately bombards the luminescent material on the anode plate to irradiate. The field emission device has a wide operating temperature range (−40° C.˜80° C.), short corresponding time (<1 ms), simple structure, low energy consumption, and meets the environmental protection requirements. Furthermore, materials such as the phosphor, luminescent glass, luminescent film, etc., can be served as luminescent material in field emission device, however, they all suffer from serious problems of low luminous efficiency, thus significantly limit the application of the field emission device, especially in the application of illumination."}
{"text": "1. Field of the Invention\nThis invention relates to hangers of the type having a receptacle for an article, such as a pot for foliage, which hanger can be suspended from an upright support element, such as an upright post and/or rail.\n2. Background Art\nThere has been a growing trend towards decorating around private homes, commercial buildings, apartments, and condominiums using hanging, potted plants. This concept is particularly popular around pool and patio decks. This type of hanger comes in a variety of different styles, ranging from a simple string-type cradle for the pots to more elaborate formed wire hangers which straddle, and are supported on, an upper edge of a support element, such as an upright post and/or a horizontally extending rail.\nOne commercially successful hanger design of the latter type is show in U.S. Pat. No. 5,390,443. In that patent, various configurations of hanger are shown, each defined by formed and welded, plastic-coated wire. Each of these hangers has a frame defining the receptacle and an integrally formed bracket, defining in conjunction with the frame, an inverted, U-shaped seat, with one leg and the base of the \"U\" defined by the bracket, and the other leg of the \"U\" defined by the frame. In a display position, the base of the \"U\" bears upon an upwardly facing surface of a support element, with the legs of the \"U\" each confronting one oppositely facing vertical surface on the support element.\nThis fixed construction, while highly commercially successful, has a number of inherent drawbacks. First of all, each bracket is configured to optimally cooperate with a single predetermined width for the support element surface. Since the surface of the support element may have a nominal width of 2 inches, 4 inches, 6 inches, 8 inches, 10 inches, etc., hangers with a wide range of bracket widths must be kept on hand to meet customer needs.\nBy reason of the receptacle size being fixed, the purveyor of these article hangers may also be required to keep on hand a large volume of article hangers having frames defining a significant range of different receptacle sizes. As the number of available bracket sizes and receptacle sizes increase, the potential combinations of these features that may be desired by the consumer increase dramatically. The purveyor thus has the option of keeping on hand a wide range of different hanger configurations or special ordering on a demand basis. The latter is not practical in that these hangers are generally relatively low cost items which may be considered \"impulse\" purchase items. Consumers are unlikely to wait any significant period of time for a special order and as a consequence sales may be lost by reason of not having the desired hanger configuration in stock.\nStocking the wide range of available configurations also has a number of drawbacks. A very significant problem with this type of hanger is that while the hangers are relatively light in weight, in an assembled state, the perimeter volume thereof is quite large. Individual boxing of the hangers often produces a bulky shape that takes up valuable warehouse and display space in the stores. Additionally, these boxes become difficult and oft times expensive to ship.\nA further problem with keeping on hand so many different configurations of hanger is that stock control is complicated. A reasonable inventory may require the assigning of a relatively large number of SKU numbers.\nFrom a manufacturing standpoint, dedicated lines, or diversions from lines, may be required to make the multiple different configurations for the hangers. This may add to the ultimate cost to the end user.\nA further problem with the fixed configuration is that the consumer, after purchase, may determine that the purchased configuration is inappropriate, necessitating a return. This may only be discovered after an attempted setup. Returns are an obvious inconvenience to any business.\nA number of the above problems have been overcome by the invention disclosed in U.S. Pat. No. 5,711,502, entitled \"Article Hanger\". The invention therein concerns a variable bracket configuration which permits a universal bracket construction for a wide range of bracket capacities.\nHowever, while the bracket configuration may be varied, the receptacle size defined by the frame is fixed. Thus, assuming different receptacle capacity is desired, the purveyor employing the variable bracket configuration must nonetheless keep on hand hangers with a number of different frame configurations, thus forcing the manufacturer and purveyors to contend with the problems noted above."}
{"text": "1. Field of the Invention\nThe present invention relates to field-programmable gate arrays, and more particularly, to a low voltage differential signaling driver for field programmable gate arrays.\n2. Description of the Related Art\nA field-programmable gate array (FPGA) is an integrated circuit (IC) that includes a two-dimensional array of general-purpose logic circuits, called cells or logic blocks, whose functions are programmable. The cells are linked to one another by programmable buses. The cell types may be small multifunction circuits (or configurable functional blocks or groups) capable of realizing all Boolean functions of a few variables. The cell types are not restricted to gates. For example, configurable functional groups typically include memory cells and connection transistors that may be used to configure logic functions such as addition, subtraction, etc., inside of the FPGA. A cell may also contain one or two flip-flops. Two types of logic cells found in FPGAs are those based on multiplexers and those based on programmable read only memory (PROM) table-lookup memories. Erasable FPGAs can be reprogrammed many times. This technology is especially convenient when developing and debugging a prototype design for a new product and for small-scale manufacture.\nAlmost all integrated circuits use input/output (I/O) buffers to connect internal circuit nodes to other circuits external to the integrated circuit. These I/O buffers can be input, output or bi-directional I/O buffers. Further, each I/O buffer may be designed to meet electrical specifications dictated by industry standards such as TTL, LVTTL, LVCMOS, GTL. It is also common for circuit designers to design each I/O buffer with multiple transistors in parallel. For example, 2-4 P-type transistors may be connected in parallel to form the pull-up section of the buffer, while 2-4 N-type transistors may be connected in parallel to form the pull down section of the buffer. Designers may then decide to use some or all of the transistors as needed by the circuit application to meet performance criteria, a particular I/O standard and noise considerations.\nThe selection of the transistors connected into the circuit is usually done by masking options such as metal, vias and contacts. Moreover, some FPGAs have used similar techniques to select one or more transistors into the I/O buffer to provide slew control. A user may configure his I/O buffer to have either fast slew or slow slew by programming an appropriate antifuse element. This feature allows the user control over speed and noise that is induced into the circuit by the switching I/O buffers.\nDifferent types of FPGAs designed by various manufacturers also feature configurable I/O buffers. These FPGAs may feature highly configurable input and output buffers, which provide support for a wide variety of I/O standards. Input buffers can be configured as either a simple buffer or as a differential amplifier input. Output buffers can be configured as either a push-pull output or as an open drain output. Selection of the desired standard is done by configuration memory bits. Further, different power supplies are provided to the I/O buffer as needed by the standard.\nHence, there is a need for an I/O that has an output buffer which can function as a low voltage differential signaling driver when used together with an adjacent output buffer."}
{"text": "A conventional fan and/or compressor assembly includes rotor blades interspaced with static vane stages which occupy a duct which defines a path down which fluid being pressurised will pass. A gap is provided between the rotatable and static components to permit relative rotation. In order to optimise efficiency of the system, it is important to minimise this gap in order to limit leakage of the air being pressurised. Typically this is achieved by either designing the gap to be as small as possible and/or by the provision of a sealing member which extends across at least part of the width of the gap. Conventionally a thin metal sealing member is located in a groove in the stator vane platform. The sealing member obstructs the gas flow through the gap and is thin enough that, in the event of excessive relative axial movement of the stator vane and rotor blade, it is easily deformed without causing damage to the vane or blade.\nHowever, since the edge of the groove in the platform in which the sealing member steps down to the sealing member, it will inherently introduce a discontinuity into the gas washed surface of the platform, thus introducing a disturbance to the flow. Also, in builds where a composite vane is used to reduce weight of the system, a groove in the platform will significantly reduce the strength of the vane and may result in composite structure coming apart. It should also be noted that the location of the sealing member is critical, since it must provide a seal without interfering with the gas flow over the platform. That is to say, it must occupy a location predefined in the design of the compressor and be neither too far inset or too proud of the gas washed surfaces of adjacent platforms. In the absence of a location groove, this is achieved using dedicated jigs, fixtures and tooling to which are expensive to produce and time consuming to use.\nHence a sealing means which can be easily located in a desired position on vane platforms is highly desirable."}
{"text": "The present embodiments relate to the operation of a hardware accelerator. More specifically, the embodiments relate to hardware accelerator memory address translation fault resolution.\nHardware accelerators are often included in processor-based systems such as computer systems to perform specific operations efficiently in hardware rather than in software. Traditionally, a hardware accelerator performs complex parallel transformations on input data, which can enhance performance of a computer based system. Additionally, in some cases, the hardware acceleration can be more power-efficient than performing the same tasks in software. Power efficiency can be even greater if the hardware accelerators are incorporated on the same semiconductor substrate (“on-chip”) as the processors. Particularly, integrating hardware accelerators onto multi-core chips such as chip multiprocessors and/or chip multithreaded processors can be efficient, because the accelerator can be shared among the cores/threads.\nTypically, a privileged layer of software in the computer system manages access to the hardware accelerator. The access management enables the hardware to be shared in a distributed manner so that various threads/cores have the opportunity to take advantage of the hardware accelerator. Additionally, the access management enables the hardware accelerator to be shared in a secure manner (e.g. preventing one thread/core from disrupting, and/or corrupting, the operation issued by another thread/core to the hardware accelerator)."}
{"text": "It is common today for people to collect billed caps, such as baseball caps and other similar novelty caps, as a hobby. Upon collecting such caps, the desire arises to display these caps in an orderly and attractive fashion. Additionally, retail stores and the like who cater to collectors and others will also desire the ability to display the different hats that they sell in an attractive fashion which will clearly display any emblem or lettering printed on the hats. Accordingly, a need has arisen for a way to support caps in an orderly and attractive fashion which is still relatively inexpensive.\nOne problem which arises when these types of hats are merely laid on a shelf or put on a basic display rack is that they are easily knocked off of the shelf or display by passers by. Therefore, there is also a need to have some type of a simple mechanism by which the hats can be secured to the support, to avoid having them easily fall off, and yet still be easily placed on and removed from the support. Further, the support should provide a simple and inexpensive way to accomplish these requirements, by being simple to set up, assemble and mount to wherever the hats will be displayed. Another problem with merely storing the caps on a shelf is that some of the caps may end up being stored in a crushed shape. Caps stored in this shape may not retain their original proper shape.\nConsequently, the need for a new and improved wall mountable hat rack for caps which is easy to use as well as inexpensive and also securely supports the caps while maintaining their proper shape is needed for collectors as well as retailers and others."}
{"text": "1. Technical Field\nThis invention relates generally to pistons for diesel engine applications, and more particularly to those having a closed oil gallery.\n2. Related Art\nIt is known in diesel engine applications to provide a piston whose piston body is formed with a closed gallery for cooling oil. The oil circulates through the gallery and cools parts of the piston which are susceptible to damage from the heat of combustion. Such cooling galleries are generally annular or ring-shaped and are provided just inside of the ring belt adjacent the top wall of the piston body. The gallery is bounded by the inner wall and closed at the bottom by a bottom wall. One or more inlets are provided for receiving cooling oil into the gallery. Various outlets are also provided for directing portions of the oil from the gallery into other regions of the piston for cooling or lubrication purposes. It is know to provide an oil passage leading from the cooling gallery to the pin bores of the pin bosses for lubricating the wrist pin. The entry point for the lubrication is rather localized and provides a limited supply of oil for lubrication which may not be present during the entire stroke of the piston.\nIt is an object of the present invention to improve the lubrication of the pin bore and closed gallery diesel piston applications.\nA piston for diesel engines constructed according to a presently preferred embodiment of the invention includes a piston body having an oil gallery and a pair of pin bosses having pin bores disposed about a common pin bore axis and pin bore surfaces which are spaced from one another about their circumference by an intervening space between the pin bosses. According to the invention, at least one recess is formed in at least one of the pin bore surfaces and extends axially across the at least one pin bore surface. At least one associated oil passage extends from the gallery to the recess for feeding oil from the gallery to the recess.\nThe invention has the advantage of providing improved lubrication to the pin bosses of diesel engine pistons having closed oil cooling galleries which serves as a reservoir for lubricating oil for the pin bosses.\nThe invention has the further advantage of providing an axial relief or recess in the pin bore surface which itself acts as a localized oil reservoir for lubricating oil adjacent the pin bore surface. The recess, fed by the oil passage, captures or wicks the oil and then feeds the oil to the pin bore surface during the full stroke of the piston, including at the top of the stroke when less oil is available in the gallery to be fed to the pin bores. At such times, the pocket of oil retained in the pin bore recess is available to feed the pin bore surfaces with lubricating oil.\nAnother advantage of the invention is that it provides a simple solution for greatly improving the lubrication of the pin bore of closed gallery diesel pistons. The recess may be formed in the pin bore surface by a simple machining operation and could be readily adapted to existing manufacturing processes during the machining of the pin bore surfaces.\nThe invention has the further advantage of prolonging the life of pistons, connecting rods and wrist pin by improving the lubrication to their interactive surfaces during operation."}
{"text": "Each day the need for economically reasonable fuel sources becomes increasingly important. Similarly the need to dispose of waste in an economical yet environmentally friendly method is also highly desirable. The present disclosure provides processes for the conversion of carbon-based material such as waste material into fuel, fuel that may be used to power turbines and/or vehicles."}
{"text": "The ability to have a desired pattern of a fluid, for example an adhesive, on a substrate is a necessity for manufacturers engaged in a variety of industries, including the packaging and plastics industries. A typical process to obtain a pattern of a fluid on a substrate is found in the packaging industry, and employs a dispensing device in the form of an adhesive gun to apply adhesive onto a substrate being moved past the dispensing device, for example, by a conveyor. A fluid dispensing system monitors conveyor movement and movement of the substrate to provide the adhesive gun with dispensing command signals.\nThe quality of the adhesive dispensing process is a subject to many variables that include general environmental conditions, the physical state of the adhesive being dispensed, the physical condition of the dispensing device, and the stability of other system parameters. Changes in such variables often result in variations in the operation of the dispensing device that in turn can produce variations from the desired pattern of the adhesive fluid on the substrate.\nThere are known devices for detecting and monitoring the placement of the pattern of a liquid, such as an adhesive dispersed by a fluid dispensing system onto a moving substrate. Some systems attempt to monitor the presence and location of adhesive on the substrate by measuring delays between dispensing gun activation and deactivation signals and signals representing leading and trailing edges of the applied adhesive, calculating a correlation between those signals, and evaluating the correlation through a statistical processing method. While such systems can effectively monitor the quality of a pattern of a fluid, those systems are not capable of monitoring a pattern of a fluid without utilizing dispensing gun activation and deactivation signals generated by a fluid dispensing system. Thus, they require a sophisticated interface with the dispensing system. Such systems are further limited to fluid dispensing systems employing an adhesive gun and are not applicable to fluid dispensing systems using other types of dispensing devices, such as a gum box.\nOther systems for testing or monitoring the quality of the pattern of an applied fluid sense an edge of an adhesive bead within a programmed window within which the edge of the bead is predicted to occur. Some of such systems require that the bead adhesive pattern that is programmed into the fluid dispersing system also be programmed into the monitoring system, while the motion of the conveyor for the substrate to which the fluid pattern is applied be monitored by a conveyor motion sensor coupled to the conveyor. Thus, these systems require both a highly skilled technician to perform the programming and the presence of a conveyor motion sensor.\nSome other systems acquire a sample fluid pattern by sampling a pattern of the fluid, the quality of which is satisfactory and then evaluating subsequent patterns by comparing time delays for leading and trailing edges of the beads of the sampled pattern with corresponding time delays extracted from the subsequent pattern. While such systems simplify the acquisition of the reference pattern data, the use of time delays as a measure for sensing the edges of the beads for comparison with the reference pattern limits the application of such systems to the monitoring of patterns of a fluid on substrates that travel with a constant line speed with respect to the dispensing and sensing devices, and requires the selection of new reference data and the associated recalibration of the system with every speed change.\nThere is accordingly a need for a fluid pattern monitoring system that effectively and reliably detects and monitors the quality of a pattern of a fluid applied by a fluid dispensing system that is applicable to a variety of dispensing devices applying the fluid to substrates moving with variable speeds with respect to the fluid dispensing device, and which operates independently from the fluid dispensing system, and is relatively easy and efficient in setup, calibration and maintenance."}
{"text": "In general, road wheels for vehicles such as automobiles are manufactured by forming and welding a steel plate material. The steel plate therefore needs a forming property, a burning property, and a weldability.\nA steel plate material of a strength 38-60 kg/mm2 grade has been used as a material of a conventional road wheel. The steel plate material of a strength 38-60 kg/mm2 grade is manufactured into road wheels by a press forming operation. The steel plate material then undergoes low temperature heat treatment to obtain a precipitating hardening effect and soft-nitriding treatment to improve the strength, respectively.\nHowever, there is a need in the automobile industry to reduce weight of components. If a road wheel is made of a material having a low strength, and the vehicle including the road wheel get a shock while traveling such as might occur by passing over a small object or a large pot-hole, the road wheel is easily distorted. Thereafter, airs in tires leaks quickly and often so fast the depletion can be termed a blowout, and reliable operation of the wheel is lost. This is a significant safety concern.\nThe strength of the material of a conventional road wheel is low so that, in order to manufacture a road wheel having a sufficient strength, the width of the material should be thick. The road wheel accordingly becomes heavy.\nWheels are conventionally manufactured by a process that includes press forming, or press molding, a metal plate. When performing a press forming of a road wheel rim using the conventional road wheel material having a low strength, with the resultant high thickness, a crack is easily generated around a cooling hole and a center hole. This reduces the percentage of manufactured tires that pass inspection and, if the flaw is not noticed, this reduces the life-span of a road wheel due to fatigue fracture propagation during the use of the road wheel.\nHowever, when a conventional high strength material, that is, of a strength higher than 38-60 kg/mm2 grade, is used for the road wheel, the press forming property is not adequate and a spring back occurs. Also, a defect such as cracking or breaking during the press forming is still generated so that it is difficult to manufacture the road wheel.\nA road wheel composition, and a method for manufacturing a road wheel using the same, which is capable of improving the strength of a finally completed road wheel while paying attention about the press forming property to form a lightweight and a durable road wheel have been greatly demanded."}
{"text": "1. Field of the Invention\nThis invention relates to a compressor member such as a scroll member for scroll compressor, to a manufacturing method for the material, and to a scroll compressor.\n2. Description of Related Art\nScroll compressors are one type of compressor used in freezing systems or air conditioning systems and are compressors performing compression by means of a stationary scroll and an orbiting scroll that orbits according to rotation of a drive shaft. Any scroll of conventional scroll compressors is made from some iron or aluminum material. Those metal scrolls, however, had a problem causing high costs for production because they require a relatively high precision to be fabricated. Furthermore, the metal scrolls are heavy in weight and therefore raise a problem that the compressors suffer from large energy losses when operating.\nPrevious inventors have developed scrolls mainly made of a resin suitable to be fabricated and light in weight in comparison with metal scrolls and have proposed various scrolls (e.g., Japanese Unexamined Patent Publication (KOKAI) No. Showa 62-199,981, Japanese Unexamined Patent Publication (KOKAI) No. Heisei 2-112,685, Japanese Unexamined Patent Publication (KOKAI) No. Heisei 2-112,688, Japanese Patent Publication (KOKOKU) No. Heisei 6-33,780).\nThose scrolls mainly made of a resin have not yet been practically produced. Although various searches for scroll materials have focused on their resistance against abrasion, research done by this inventor discovered that there were other practical characteristics to be considered in addition to resistance against abrasion. That is, the scrolls in the scroll compressors are in contact with refrigerant and refrigerating machine oil at a high temperature and a high pressure during operation of the compressors. Scrolls mainly made of a resin in contact with the refrigerant and refrigerating machine oil at the high temperature and the high pressure change in dimension, thereby causing a problem in that the scroll prevents the compressor from operating efficiently and possibly makes it inoperable. The research concluded that to make the scroll practical, the scroll needs to be made of a material having dimensional and chemical stability under high temperature and high pressure refrigerant and refrigerating machine oil, as well as abrasion resistance.\nVarious engineering plastics are considered to be used for some mechanical parts other than scrolls in scroll compressors, and are in fact, used in many fields. Particularly, those are frequently parts of automobiles. For example, known as mainly made of a phenol aralkyl resin are pulleys (e.g., Japanese Patent Publication (KOKOKU) No. Heisei 4-38,944, Japanese Patent Publication (KOKOKU) No. Heisei 4-38,945), rotors for pump (e.g., Publications of U.S. Pat. Nos. 2,643,476 and 2,643,477), and impellers for pump (e.g., Japanese Unexamined Patent Publication (KOKAI) No. Heisei 8-93,690, Japanese Unexamined Patent Publication (KOKAI) No. Heisei 9-112,489).\nHowever, in all of the above cases, no material is selected in consideration for dimensional and chemical stability under high temperature and high pressure refrigerant and refrigerating machine oil, as well as abrasion resistance."}
{"text": "1. Technical Field\nThis disclosure relates to semiconductor fabrication and more particularly, to a method for applying and removing an inorganic anti-reflection coating for semiconductors.\n2. Description of the Related Art\nAnti-reflection coatings (ARCs) are typically deposited in semiconductor fabrication prior to the deposition of a resist material. ARC coatings absorb radiation to form an optical opaque film which enhances the optical contrast of the imaging resist. ARC coatings effectively reduce reflection (anti-reflection) of the incident radiation back into the overlying resist layer. This prevents overexposure of the resist material.\nDielectric anti-reflection coatings (DARCs) have become more attractive for deep ultraviolet (DUV) lithographic ARC applications. It may be desirable for DARCs to replace organic ARC processes in semiconductor industry depending on the application of the coating. DARC processes have many advantages over organic ARC processes. Both the thickness and the chemical composition of the DARC can be optimized to minimize optical reflections from an underlying film stack to provide lithographic control of feature sizes. DARC can provide higher dry etch selectivity to DUV resist in mask open processes for at least two reasons. One, the DARC layer may be substantially thinner than organic ARC. Second, the etch selectivity of DARC to resist can be optimized to be greater than one, whereas organic ARC selectivity to resists is usually limited to less than or equal to one. Additionally, since the DARC thickness can be precisely controlled, DARC, DARC have the potential for better critical dimension control during the ARC open processes.\nHowever, DARC is not usable for some applications because it is difficult to remove DARC after a lithographic process. DARC is more difficult to remove relative to removing an organic ARC or a resist. For example, when DARC is used on the top of a silicon nitride layer such as for processing at an active area (AA) level of a semiconductor device, it is extremely difficult to eliminate DARC without nitride loss because DARC is very similar to nitride in terms of its chemical composition.\nOrganic ARC layers tend to be thick since they are etched at about a same rate as the resist layer deposited thereon. Typically, organic ARC layers tend to be about 900 xcex94 to about 1100 xcex94 in thickness. This thickness is not desirable since the thickness may adversely affect dimension control for patterning structures in the semiconductor device.\nTherefore, a need exists for a for a method for forming a DARC layer and removing the DARC layer in semiconductor fabrication processes. A further need exists for a thinner anti-reflection layer with higher selectivity to resist.\nIn accordance with the present invention, a method for employing and removing inorganic anti-reflection coatings, includes the steps of providing a first dielectric layer on a semiconductor device structure to be processed, the first dielectric layer being selectively removable relative to the semiconductor device structure, and forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer, the DARC being selectively removable relative to the first dielectric layer. A resist layer is patterned on the DARC. The resist is selectively removable relative to the DARC. The semiconductor device structure is etched, and the resist layer, the DARC and the first dielectric layer are selectively removed.\nA method for etching an active area, in accordance with the invention includes providing a substrate having deep trenches formed therein and a pad stack formed thereon, the trenches including storage nodes and buried straps formed in the trenches. A first dielectric layer is provided on the pad stack, on sidewalls of the trenches and over the buried straps. The first dielectric layer is selectively removable relative to the pad stack, the sidewalls of the trenches and the buried straps. An inorganic dielectric anti-reflection coating (DARC) is formed on the first dielectric layer. The DARC is selectively removable relative to the first dielectric layer. A resist layer is patterned on the DARC, and the resist is selectively removable relative to the DARC. Positions are etched for shallow trench isolation regions in accordance with the patterned resist layer. The resist layer, the DARC and the first dielectric layer are selectively removed.\nIn other methods, the step of providing a first dielectric layer may include the step of forming a conformal first dielectric layer. The method may include the step of forming an overhang structure with the first dielectric layer to prevent at least some coverage of vertical surfaces by the DARC. The first dielectric layer may include a silicon oxide or a polymer. The step of forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer may include the step of forming the DARC including a silicon oxynitride. The step of forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer may include the step of annealing the DARC to prevent interactions between the resist layer and the DARC. The step of forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer may also include the step of depositing the DARC by one of a chemical vapor deposition process and a physical vapor deposition process. The method may include the step of forming a second dielectric layer on the DARC, the second dielectric layer being selectively removable relative to the DARC.\nA method for employing inorganic anti-reflection coatings, in accordance with the present invention, includes providing a first dielectric layer on a semiconductor device structure to be processed, the first dielectric layer being selectively removable relative to the semiconductor device structure, and forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer, the DARC being selectively removable relative to the first dielectric layer. A resist layer is patterned on the DARC, the resist being selectively removable relative to the DARC. A portion of the DARC and the first dielectric layer are etched in accordance with resist layer. The resist layer is selectively removed with respect to the DARC, and the semiconductor device structure is etched by employing the DARC as a hard mask.\nIn other methods, the step of selectively removing the DARC and the first dielectric layer is included. The step of etching the semiconductor device structure by employing the DARC as a hard mask may include the step of consuming the DARC during the etching of the semiconductor device structure. The first dielectric layer may include a silicon oxide. The step of forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer may include the step of annealing the DARC to prevent interactions between the resist layer and the DARC. The step of forming an inorganic dielectric anti-reflection coating (DARC) on the first dielectric layer may include the step of forming the DARC with silicon oxynitride.\nThese and other objects, features and advantages of the present invention will become apparent from the following detailed description of illustrative embodiments thereof, which is to be read in connection with the accompanying drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of controlling a clutch pressure of a transmission structured so as to engage a hydraulic clutch so as to set a stage for shifting a speed, and a control apparatus therefor.\n2. Description of the Related Art\nThe transmission generally has a main transmission which is provided with a clutch for a plurality of stages for shifting the speed including a backward movement, and a sub transmission which is provided with a clutch for a high speed and a low speed, and there has been known a structure which performs a speed change in forward and backward directions by a combination between the main transmission and the sub transmission and a structure which is provided with a clutch for a forward movement stage for shifting the speed and a clutch for a backward movement stage for shifting a speed.\nIn recent years, a large-size transmission has been developed together with an enlargement of a vehicle.\nThe applicant of the present application has filed Japanese Patent Publication No. 5-65731 and PCT International Publication No. WO98/01687 with respect to the following technique, as a clutch hydraulic pressure control apparatus in a transmission corresponding to an enlargement of the vehicle.\nThere has been disclosed a method which is provided with an electromagnetic proportional pressure control valve and a flow amount detecting valve for each of speed change clutches of a clutch hydraulic pressure control apparatus, the method detecting a filling state of a clutch operating fluid within a piston volume of the clutch by the flow amount detecting valve, operating the electromagnetic proportional pressure control valve in real time at the same time of being filled with fluid so as to perform a gradual increase of a fluid pressure, performing a clutch engagement control, and measuring a time for being filled with the fluid at the preceding clutch engagement so that a timing for the clutch engagement does not become different even when the time for filling a clutch chamber is changed due to an abrasion of a clutch disc, a dispersion of a part processing accuracy and a variation of an engine speed so as to automatically set to a suitable target filling time after comparing with the suitable target filling time.\nHowever, because in the technique described in Japanese Patent Publication No. 5-65731 and PCT International Publication No. WO98/01687, the fluid is caused to flow into the clutch chamber through throttling passages of the electromagnetic proportional pressure control valve and the flow amount detecting valve, the filling time of the fluid into the clutch chamber is elongated and the clutch engagement is delayed when it is applied to a clutch of a large capacity piston volume in a large-size transmission, so that an engagement time lag is generated. In particular, in the case of a construction vehicle, it is necessary to frequently repeat a speed change between a forward movement first speed and a forward movement second speed for excavating earth and sand. Accordingly, it is necessary that an amount of fluid to be supplied to the clutch chamber is set to a large amount when starting the speed change and to a small amount immediately before the clutch is engaged, thereby performing a rapid speed change and a reduction of the speed change shock.\nFurther, because the operation is performed while frequently repeating an engagement and a disengagement of the clutch, the oil is again started to be supplied to the clutch before the discharge of the oil is competed in the case that a discharging time of the oil within the clutch chamber at the disengaging time is long, so that a dispersion is generated in the time for the clutch engagement in the same clutch, an there is a problem that a speed change shock is easily generated because a lot of oil is supplied to the clutch chamber in which the oil is still left.\nIt is possible to supply a lot of oil at a time of starting the speed change and make the discharging time of the soil short at a time of disengaging the clutch by enlarging a capacity of the electromagnetic proportional pressure control valve and the flow amount detecting valve. However, a large space is required and costs are high."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nIn some instances, it may be desirable to fix a fractured bone with a bone plate. In many examples, bone plates are configured with threaded apertures having different diameters that may be used to accept bone screws having different diameters and/or lengths. As can be appreciated, it is desirable to incorporate a bone screw having the correct length such that a distal tip of the bone screw does not extend beyond the far side of the bone. Similarly, implanting a bone screw that has a length that is too short can reduce the effective purchase of the bone screw, which is also undesirable.\nIn this regard, optimization of the length of a bone screw can be critical for achieving proper fixation of a bone plate. A depth gage can be used to measure the depth of a prepared bone hole to determine the appropriate size bone screw. Often, each size of bone screw requires a specific depth gage. In this regard, a surgeon may be required to utilize different depth gages for measuring the depth of a prepared bone hole relative to a given threaded aperture of the bone plate. It is desirable to quickly and efficiently measure the depth of a hole prepared in bone that may be configured to receive a bone screw having a first diameter as well as another hole prepared in the bone that may be configured to receive a bone screw having a second distinct diameter."}
{"text": "It is known from e.g. U.S. Pat. No. 4,025,575 that light oxygenates can be converted to the lower olefins ethylene and propylene using H-ZSM-5 zeolite as a catalyst.\nU.S. Pat. No. 3,911,041 describes that methanol and DME can be converted to a reaction product containing olefins using a zeolite catalyst containing at least 0.78 wt-% phosphorus incorporated with the crystal structure of the zeolite. The zeolite used in the process of U.S. Pat. No. 3,911,041 may be ZSM-5. U.S. Pat. No. 3,911,041 further teaches that the activity of the phosphorus-containing zeolite catalyst can be increased by depositing zinc (Zn) on the zeolite.\nU.S. Pat. No. 5,367,100 describes a method for the conversion of methanol or dimethyl ether to light olefins using a zeolite ZSM-5 based catalyst containing at least 0.7 wt-% phosphorus and at least 0.97 wt-% rare earth elements which are incorporated within the structure of the catalyst. The rare earth elements are preferably rich in lanthanum so that the content of lanthanum in the catalyst preferably is between 2.5 and 3.5 wt-%.\nU.S. Pat. No. 4,049,573 describes a number of different boron or magnesium-comprising catalyst compositions useful for converting monohydric alcohols and their ethers to a hydrocarbon mixture ring in C2-C3 olefins and mononuclear aromatics.\nThe oxygenate-to-olefin catalysts of the prior art have the disadvantage that they have a relatively low selectivity for C2-C3 olefins and/or become quickly deactivated by coke deposition on the catalyst surface."}
{"text": "1. Field of the Invention\nThis invention relates generally to audio enhancement and more particularly to systems and methods to enhance audio signals to compensate for phase distortion.\n2. Description of the Related Art\nWhen voice or music is reproduced through transducers, such as loudspeakers, which exhibit both spring and mass in their construction, the transducers produce phase shifting of the audio signal throughout the frequency spectrum within which the transducers operate.\nIdeally, the cone of a loudspeaker operates as a unit. At a frequency determined by the design of the loudspeaker, the cone propagates a radial wave outward from the voice coil region to the surround region. The loudspeaker delays higher frequency waves with respect to lower frequency waves. In an embodiment, the amount of delay is dependent on the speed of sound permitted by the speaker material. The delay produces phase shifting throughout the frequency spectrum within which the transducers operate.\nThis has been realized by loudspeaker manufacturers, but little can be done in the manufacturing process to alleviate this problem. Components of speakers, such as cones and diaphragms, have mass. Speaker centering and restoring devices, such as spiders and surrounds, exhibit spring constants. Further, air itself exhibits a spring constant. This combination produces phase shifting throughout the frequency range.\nThe phase shifting produces significant changes in the waveform, which the transducer is attempting to reproduce. The result is a phase distortion of the audio signal, which is distinctly audible."}
{"text": "1. Field of the Invention\nThe present invention relates to the design of Ethernet passive optical networks. More specifically, the present invention relates to a method and apparatus for facilitating asymmetric line rates in an Ethernet passive optical network.\n2. Related Art\nIn order to keep pace with increasing Internet traffic, optical fibers and associated optical transmission equipment have been widely deployed, substantially increasing the capacity of backbone networks. This increase in the capacity of backbone networks, however, has not been matched by a corresponding increase in the capacity of access networks. Even with broadband solutions, such as digital subscriber line (DSL) and cable modem (CM), the limited bandwidth offered by current access networks creates a severe bottleneck in delivering high bandwidth to end users.\nAmong the different technologies that are presently being developed, Ethernet passive optical networks (EPONs) are one of the best candidates for next-generation access networks. EPONs combine ubiquitous Ethernet technology with inexpensive passive optics. Hence, they offer the simplicity and scalability of Ethernet with the cost-efficiency and high capacity of passive optics. In particular, due to the high bandwidth of optical fibers, EPONs are capable of accommodating broadband voice, data, and video traffic simultaneously. Such integrated service is difficult to provide with DSL or CM technology. Furthermore, EPONs are more suitable for Internet Protocol (IP) traffic, because Ethernet frames can directly encapsulate native IP packets with different sizes, whereas ATM passive optical networks (APONs) use fixed-size ATM cells and consequently require packet fragmentation and reassembly.\nTypically, EPONs are used in the “first mile” of the network, which provides connectivity between the service provider's central offices and business or residential subscribers. Logically, the first mile is a point-to-multipoint network, with a central office servicing a number of subscribers. A tree topology can be used in an EPON, wherein one fiber couples the central office to a passive optical splitter/combiner. This passive optical splitter/combiner divides and distributes downstream optical signals to subscribers and combines upstream optical signals from subscribers (see FIG. 1).\nTransmissions within an EPON are typically performed between an optical line terminal (OLT) and optical network units (ONUs) (see FIG. 2). The OLT generally resides in the central office and couples the optical access network to a metro backbone, which is typically an external network belonging to an Internet Service Provider (ISP) or a local exchange carrier. An ONU can be located either at the curb or at an end-user location, and can provide broadband voice, data, and video services. ONUs are typically coupled to a one-by-N (1×N) passive optical coupler, where N is the number of ONUs, and the passive optical coupler is typically coupled to the OLT through a single optical link. (Note that one may use a number of cascaded optical splitters/couplers.) This configuration can significantly save the number of fibers and amount of hardware required by EPONs.\nCommunications within an EPON include downstream traffic (from OLT to ONUs) and upstream traffic (from ONUs to OLT). In the downstream direction, because of the broadcast nature of the 1×N passive optical coupler, data frames are broadcast by the OLT to all ONUs and are subsequently selectively extracted by their destination ONUs. In the upstream direction, the ONUs need to share channel capacity and resources, because there is only one link coupling the passive optical coupler with the OLT.\nTo interoperate with other Ethernet equipment, an EPON ideally complies with the IEEE 802 standards. Accordingly, the EPON architecture is standardized in the IEEE 802.3ah standard. This standard only provides for symmetric line rates, i.e., both upstream and downstream communications are performed at the 1.25 Gbps line rate. EPON users, however, often desire different upstream and downstream line rates. Unfortunately, there is no existing EPON architecture that allows asymmetric line rates.\nHence, what is needed is a method and an apparatus for facilitating asymmetric line rates in an EPON which allows a service provider to provide more diversified services."}
{"text": "The present invention relates to a bag for an infusion solution particularly having a structure capable of easily mixing, at a time of use, an infusion solution with another infusion solution or a medicament other than the infusion solution which have been preliminarily independently sealed in the bag and also relates to a method of manufacturing such infusion solution bag.\nIn a case where it is required to prescribe an infusion solution, it has often be performed to prescribe a solution which has been prepared by mixing more than two kinds of infusion solutions or to prescribe a solution which was prepared by mixing a medicament solution other than the infusion solution with the infusion solution or dissolving powder of a medicament into the infusion solution. However, when the solution prepared by mixing the infusion solution with another infusion solution or medicament and the mixed solution is preserved, there is a fear such that the mixed solution may be denatured or discolored. Because of this reason, it is necessary to preliminarily preserve these solutions or medicaments as they are and mix them just before the prescription thereof. For example, with a dialysis solution utilized for peritoneal dialysis, it is difficult to preserve it for a long time if a pH thereof be not adjusted to weak acid, but it gives strong stimulus to a human organism in a case where it has not the neutral pH. Accordingly, it is necessary to add an alkali solution to the dialysis solution just before the prescription to make the pH thereof neutral. For example, in a case where a mixture solution of glucose infusion solution and amino acid infusion solution has been left as it is for a long time, Maillard reaction will be caused and the mixture solution will be discolored. In this example, it is also necessary to mix these solutions just before the prescription.\nIn view of the above matters, there has been provided an infusion solution bag having an inner space which is divided into a plurality of small chambers or sections into which infusion solutions, or infusion solution and a medicament other than the infusion solution, are accommodated respectively. These small chambers are sectioned by boundary portions which are sealed to be openable. These infusion solutions or medicament solutions stored in the respective chambers can be mixed together by opening the boundary portions just before the use.\nIn one type infusion solution bag of conventional structure, a shell portion of the infusion solution bag is pinched from the outer side thereof at several portions by using pinching members to divide the inner space of the bag into a plurality of small chambers or sections and various infusion solutions or the like are accommodated into these divided sections respectively. The pinching members are removed at the time of prescription to mix these infusion solutions together.\nIn another type of conventional structure, opposing film materials constituting both surfaces of an infusion solution bag are closed, through a heat sealing process, at boundary portions by which the inner space of the infusion solution bag is divided into a plurality of small sections. A port forming member is mounted to a part of each boundary portion, infusion solutions or other medicament solutions are accommodated in the small sections respectively, and these solutions are mixed at the infusion solution prescription time by breaking, for example, bending, the port forming member.\nHowever, according to the conventional structures mentioned above, it is necessary to preliminarily prepare pinching members or port forming members and attach them to the infusion solution bag, resulting in increasing of troublesome working and increasing of manufacturing cost. Furthermore, it is difficult to strictly section the respective small chambers in desired shapes, and moreover, in the case of using the port forming member, since the thus formed port has a small opened area, it takes much time to complete the mixing of the infusion solutions accommodated in the respective sections.\nStill furthermore, prior art provides a method in which opposing film materials, constituting both surfaces of an infusion solution bag, are closed, through a heat sealing process, with a resin film being interposed therebetween for weakening a sealing strength, at boundary portions by which the inner space of the infusion solution bag is divided into a plurality of small sections. Infusion solutions or other medicament solutions are accommodated into these small sections or chambers, and when it is required to mix these solutions in the respective small sections together, the weakly sealed part of the boundary portion is separated by applying pressure to the bag body and, then, the solutions in the respective sections are mixed through the opened boundary portion.\nHowever, in this conventional structure, it is also necessary to prepare a resin film formed of a relatively specific material having a weakened sealing strength and dispose such resin film to a desired position of the boundary portion sectioning the respective small chambers of the bag, also involving troublesome working and increasing the manufacturing cost. Furthermore, since it is difficult to arrange the resin film having the weakened strength in a desired pattern, it is also difficult to section the small chambers so as to provide desired shapes. Still furthermore, in this structure, it is necessary to select the material for the resin film for the adjustment of the sealing strength of the boundary portion to be separated and to suitably adjust the heating temperature and pressure at the sealing time. However, it is extremely difficult to finely adjust the sealing strength of the resin film by selecting the material therefor, and moreover, in the heat sealing process, it is difficult to stably control the sealing strength by adjusting the heating temperature and pressure, which may result in that the thus manufactured infusion solution bag has no uniform sealed strength at the boundary portions to be separated."}
{"text": "The invention relates to a fluid control valve; particularly to an anti-backflow, anti-siphon vacuum breaker valve.\nThere are many fluid flow systems that supply fluid from a clean source for use in a potentially contaminated fluid environment. These systems are regulated by local laws to comply with certain industrial codes, standards and specifications to protect the clean source of fluid for other uses and users. These codes and laws generally require that none of the contaminated fluid is permitted to backflow into the clean (or potable) source of fluid. This requires a fluid control valve having a vacuum breaker vent to prevent any back-siphoning of the contaminated fluid, as well as a direct anti-backflow valve to prevent any backflow of the contaminated fluid from flowing into the clean source of fluid.\nSome examples of such systems include agricultural water irrigation systems, home dialysis installations and other medical applications, solvent cleaning systems, industrial water and chemical applications, and numerous home and industrial plumbing and sewage handling applications.\nThere are many prior art devices, that have been developed for use in these applications and environments. Some examples of such devices include the Watts 188A for irrigation systems, the Watts 288A for hot or cold water anti-siphon vacuum breakers, and the Watts N338 for hot and cold water anti-siphon vacuum breaker having bottom inlet and outlet connections. There are numerous other manufacturers and suppliers of such devices.\nThe prior art devices, because of the venting design, usually require that the device be installed as the last component before the discharge end, and with no downstream shutoff valve, and not less than 12 inches above the discharge outlet. These requirements present design limitations for many fluid control systems.\nA common problem associated with the prior art devices is that in certain backflow conditions, even though they comply with the anti-backflow requirements to protect the clean source, they permit the contaminated backflow fluid to flow through the vent into the open environment. This may not present a problem in outside irrigation applications or in industrial controlled spillage environments; but this is a major problem in many laboratory and home installations.\nAs shown in FIG. 1 (Prior Art), a typical example of a prior art device includes an anti-siphon vacuum breaker valve 10, having an inlet 12, an outlet 14, and an atmospheric vent 16. The valve includes an internal chamber having an inlet valve seat 18 and an upper vent seat 20, and a dual action disc valve 22. The valve is slidable on a vertical guide stem 24 and has a lower valve seal 26 and an upper vent seal 28.\nIn operations, the lower valve seat 18 is normally closed, and the upper vent seat 20 is normally open; when water pressure is supplied to the inlet (not shown), the disc valve 22 is elevated on the stem to open the water inlet valve seat 18 and the upper seal 28 engages the vent seat 20 to close the vent 16. When the water pressure at the inlet is discontinued, the disc valve 22 returns to the lower position. Also in operation, when backflow is applied through the outlet 14, the valve seal 26 engages the valve seat 18 and prevents contaminated flow from entering the inlet source fluid, but the vent seat 20 can not be sealed; and the contaminated fluid is permitted to flow through the vent 16 into the environment and onto the floor of the home or laboratory. This is a major mess, and potential health threat to the home or laboratory.\nIn view of the foregoing, it is an object of the present invention to provide an anti-siphon, anti-backflow, vacuum breaker valve that prevents any backflow of contaminated fluid into a clean (potable) fluid source, and prevents any backflow of contaminated fluid through the vent and into the environment and onto the floor of the home, laboratory or commercial facility.\nIt is another object to provide an anti-siphon anti-backflow vacuum breaker valve that is compact, economical to produce, and reliable and durable in operation.\nIt is another object to provide an anti-siphon anti-backflow vacuum breaker valve that can be installed in any sequence, including forward of a shut-off valve, and at any elevation, relative to the discharge end, in a fluid flow system."}
{"text": "The present invention relates to integrated circuit (IC) design, and more specifically, to a method and system for clock tree construction, an IC and a fabrication method thereof.\nA typical IC may comprise a large amount of logic elements and other circuits for implementing IC logic functionality. Further, an IC chip may comprise a clock tree (i.e., a clock signal distribution network) for distributing a clock signal received at an input to all clock sinks that are “clocked” by the clock signal. A clock tree may comprise wires, buffers etc, to distribute the “clock signal” that controls the timing and operation of logical elements and other circuits of the IC. Clock sinks (or sinks) refer to logic elements or other circuits such as registers (flip-flops), RAM and latches, controlled by a clock signal, such that they add capacitance to the clock tree. Those sinks can change their states in response to clock signal pulses, and the IC synchronizes state changes in various sinks in a clock domain by clocking them with the same clock signal.\nClock skew is a significant aspect in assessing clock tree performance and quality. Clock skew generally refers to the difference (delay) between arrival times at any two clock sinks of a clock signal from an external clock source. Due to different path lengths of various branches to the respective clock sinks in a clock tree, there may be some clock skews between those various clock sinks. Further, in order to deliver the clock signal to every region of an IC, clock cells (buffers, for example) are generally inserted in the clock tree to amplify and/or retransmit the clock signal. However, because each clock cell has an intrinsic delay, it may cause a certain clock skew also. Thus, controlling or restricting the level number of buffers in a clock tree is one way for improving clock tree performance and IC design quality. Theoretically, smaller clock skews can be obtained, if there are less but the same levels of buffers contained in each branch leading to various clock sinks in a clock tree. However, the above assumption cannot be satisfied in practical IC design in many situations. With the technical evolution of digital IC design, a common path method is becoming more important for improving clock tree's skew and timing. A common path generally refers to a path consisted of buffers that are shared by multiple sinks in a clock tree. The longer the common path is, the smaller the clock skew of a clock signal arriving at sinks is. Traditional techniques employ a method to maximize the common path, that is, to allow sinks in a clock tree to share buffers at various levels as much as possible. In principle, the more buffers shared in a clock tree, the longer the common path will be. As a result, the performance of the clock tree may be optimized, and the quality of the IC designed may be improved. Other means for clock tree optimization include, for example, utilizing high performance clock cells (elements) capable of reducing clock skew, and the like."}
{"text": "Stems are used in prosthetic joint implants to anchor the implant in a bone cavity. The bone is typically prepared to receive the stem by milling, drilling, or reaming the bone to form a cavity sixed and shaped for receiving the stem of the implant. The stem is driven into the prepared cavity of a bone, with its proximal end extending out of the cavity and carrying a joint bearing surface, which may be part of an assembly separately attached or coupled to the proximal stem. Many stems are fluted, i.e., they have distally extending longitudinal splines or ridges that cut their own force fit and provide greater stability and anchoring when the stem is inserted into the prepared cavity. Other stems may rely on cement fixation for a major portion of their length, but these also generally have at least some portion in the upper region that has a muting fit to the adjacent prepared bone surface. This upper portion must be placed in a precisely defined rotational orientation in order to seat properly when the stem is fully inserted.\nTypically, once a fluted force-fitted stem is driven a certain distance into a prepared cavity, it becomes extremely difficult to rotate the stem to reorient the position of the upper portion. Instead, the stem must be withdrawn and reinserted, a process which can damage the bone and complicate, or increase the duration of, the operation. The stem may be aligned by marking the bone and the stem, and then matching or sighting along the respective marks to align the stem with respect to the bone as it is driven in. For example, to install an artificial hip prosthesis, the surgeon may measure the mechanical axis of the femur, perform a resection of the proximal femoral end, and then mark the plane of the mechanical axis on the periphery of the resected surface to indicate the required angular orientation of the ball joint component. The mark is then used to orient the stem as it is inserted. A drawback of this method is the imprecision of the alignment process. Because the stem seating surface is at the top of the stem, and the relevant bone markings lie at the distal end, the respective landmarks on the stem and bone are initially not in close proximity to each other. Thus, parallax and other problems associated with alignment by eye may result. Also, the stem may wander from its initial aligned position as it is being advanced into the prepared cavity.\nThe above-referenced patent application addresses this problem by providing prosthesis stem components having an axially-running flute or slot, or external splines, and providing a separate alignment body which nestles in the prepared cavity in a defined orientation, and keys to flute or slot such that the stem is held in the desired rotational plane as it is inserted. The alignment body or tool must be withdrawn before the stem is fully inserted, but after the stem is sufficiently well engaged to hold its own alignment. Various embodiments may control the degree of purchase of the stem and the phase where removal of the tool is effected, or may adapt those guides to different prosthesis systems. This may be done by employing a somewhat flexible polymer for the flute-engaging body, forming it with a jaw that flexes open for removal, by employing only a narrow key or guide, or by adjusting the length or diameter of the fluted region or length of slot and depth of bite to achieve dependable insertion depth and grip before tool removal is necessary.\nHowever, there remains a need to provide tools and methods for aligning the stem of an orthopedic implant during insertion into a prepared bone cavity."}
{"text": "Measurements of analytes in blood is commonly performed by sampling blood from patients and analyzing said samples in a laboratory, often situated at a location remote from the ward. E.g. for glucose analysis there are available special reagent sticks usable for measuring on site i.e. in the ward. However, the accuracy of such measurements is questionable, and the error could be 10-20% at best.\nOften it is necessary to perform several sequential measurements over periods of several hours, which is very labor intensive. Furthermore, the risk for errors because of the human intervention is evident, and the low accuracy is of course also a drawback in this regard.\nFor the purposes of this application, the term “biosensor” means any device having a portion which interacts with biological or biochemical material, and has the capability to generate a signal indicative of a change in some parameter of said biological or biochemical material as a consequence of said interaction.\nWhen analytes such as glucose, urea, lactate, ATP, glycerol, creatinine and pyruvate in biological samples, such as blood, plasma or serum are analyzed using biosensor techniques based on immobilization of enzymes, the sensor surface will be exposed to a certain amount of sample during a certain time sufficient to achieve an adequate sensor response. It is well known that the sensor response gradually will degrade because of fouling of the surface. This in its turn is a consequence of said exposure and the interaction between the surface and the substances present in the sample that occurs. The chemical and physical composition of the sample is thereby of importance, the sample i.a. comprising red cells, blood platelets, macromolecules, electrolytes, lipids, red/ox-compounds etc. It is also known that the support material for the enzyme immobilization in biosensors based on enzyme column technology is fouled by the substances present in the sample.\nIn cases where selective membranes are used for protection of the sensor surface of biosensors based on enzyme electrode technology, said membranes are also fouled by such substances.\nThis fouling influences the sensor response by substantially reducing the life and stability of the biosensor."}
{"text": "Active implantable medical devices (IMDs) may be used to treat a variety of symptoms or diseases. For example, an IMD may be used to generate therapeutic electrostimulation for delivery to one or more tissue sites, such as including one or more cardiovascular or neural targets. Such electrostimulation therapy may include one or more of a cardiac pacing therapy, a cardioversion or defibrillation therapy, a neural stimulation therapy that may include an autonomic modulation therapy (AMT), or one or more other electrostimulation therapies.\nIn generally-available IMDs, electrostimulation may be generated using power supplies or other circuitry dedicated or hardwired for use with a particular electrode combination, therapy mode, or tissue site. Such therapy modes may correspond to using a power supply circuit configured only for operation in a controlled-current mode for neurostimulation or a power supply circuit configured only for operation in a controlled-voltage mode for cardiac pacing. If multiple tissue sites are to be stimulated, such power supply or other circuitry generally includes dedicated power supplies for each lead or electrode configuration corresponding to a respective site. Such generally-available IMDs may provide only limited options regarding available electrode combinations or therapy modes. For example, an IMD configured for current-mode stimulation is generally not reconfigurable to provide voltage-mode stimulation."}
{"text": "1. Technical Field\nThe present invention generally relates to tools for developing software for international use and in particular to multi-language software development. Still more particularly, the present invention relates to an improved system and method for performing language translation in computer software.\n2. Description of the Related Art\nSoftware Internationalization\nAs computers have become more prevalent, it has become desirable for software developers to be able to market their products to those people who do not speak the native language of the software developers. In particular, it is desirable that software developed in the English language be available to those persons, both in the United States and in the rest of the world, that do not speak English. Accordingly, many software applications that are developed in English are later translated for use by non-English speakers.\nThe process of translating a software package into another (or more than one other) language is time-consuming and expensive. Each text message, menu, and button must be translated to allow the user to operate the program. The most direct way to do this is to search the entire program source code for every text string, i.e., every string of characters that would be displayed to the user, and translating each of these to the new language.\nThis approach has several problems. One problem is that the use of this method means that the software must be specifically translated and compiled for each intended language. This, of course, is an expensive process in itself, and means that any change in the source code requires each language version of the code to be edited and recompiled.\nFIG. 2 shows an example of hardcoded labels in a Java application. In this Figure, window 200 shown with two graphical user interface (GUI) \"buttons\" 210 and 230. These buttons are created using a hard-coded Java \"JButton\" call; for example, code-line 220 calls JButton with the string argument \"Ok\", creating button 210 with text label \"Ok\". Similarly, code-line 240 calls JButton with the string argument \"Cancel\", creating button 230 with text label \"Cancel\". To translate a software application which uses hard-coded strings as in FIG. 2, each line of the code must be parsed and each displayed string must be individually translated for each target language.\nOne solution to this problem is the use of separate localization files, in which the text strings that are to be displayed are stored separately from the executable code itself. As the software is executed, the text for every given display screen is simply read from the localization files, in whichever language is stored in the file. In this manner, the text in the localization file can be translated without disturbing the executable, and the executable can be changed or replaced without disturbing the translated text (except, of course, that if the text to be displayed changes, the corresponding entries in the localization files must also be changed). The localization files may be in any number of formats, including compiled message catalogs, Java resource bundles, HTML files, and many others.\nThe use of localization files provides an easy way to do multiple translations, but presents an additional problem in that the translation is performed on text isolated in a separate file, without any kind of context for the translator. Because this is the case, the translation will often contain errors which would not be present if the translator had been able to perform the translation according to the context in which the text appears.\nHowever the translation is handled, each screen of the program in operation must then be proofread to ensure that the translated text is properly displayed and properly translates the text in the context of the display. Therefore, it has become common practice to hire individuals with backgrounds in other languages to proofread each screen of the translated program to detect any display, grammar, or other translation errors. This second phase is called Translation Verification Testing (TVT).\nFIG. 3 shows a chart of a translation and testing cycle according to conventional methods. A typical Company A, located in Country A, builds the application intended for Country B (310). Next, Company A packages and ships the application to Company B, located in Country B (320). Company B will then execute the application to test for any translation or internationalization errors (330). Company B will then do TVT testing, note any translation errors, and make any necessary retranslations (340). Company B then ships the corrections back to company A, and the process repeats until no more errors are found.\nAccording to this conventional process, each time an error is corrected or a change made in translation in Country B, the product must be rebuilt and repackaged in Country A before it is retested. Because of this, the TVT process is necessarily an expensive and time-consuming process. It would therefore be desirable to provide a means to provide additional information to the translator, and contextual information in particular, so that the initial translation is as accurate as possible. By doing so, the cost and time required for TVT testing can be reduced or eliminated.\nJava\nJava is two things: a programming language and a platform. In one aspect, Java is a high-level programming language, and is unusual in that each Java program is both compiled and interpreted, to make it accessible to many platforms. Java is also a software-only platform that runs on top of other, hardware-based platforms. The Java platform has two components; the Java Virtual Machine (Java VM), and the Java Application Programming Interface (Java API).\nThe Java VM is the base for the Java platform and is ported onto various hardware-based platforms. The Java API is a large collection of ready-made software components that provide many useful capabilities, such as graphical user interface (GUI) widgets. The Java API is grouped into libraries (packages) of related components.\nJavaBeans brings component technology to the Java platform. JavaBeans API supports reuseable, platform-independent components. Using JavaBeans-compliant application builder tools, components can be combined into applets, applications, or composite components. JavaBean components are known as Beans."}
{"text": "1. Field of the Invention\nThe present invention generally relates to the winding of tape around reels in a tape drive and, more particularly, to an integrated reel hub and motor shaft assembly and fabrication method therefore that reduce radial and axial run-out of the tape relative to the reel hub to limit position errors and head-to-tape spacing complications.\n2. Relevant Background\nTape drives have been widely employed in industry for over thirty years due to their ability to store large amounts of data on a relatively small and inexpensive removable format. Typically, tape drives use a storage tape having a thin film of magnetic material which is wound between a pair of tape reels as data is transferred to or from the tape media via a read/write tape head assembly. In one arrangement, one of the reels (e.g., the “take-up” reel) is part of the tape drive while the other reel (e.g., the “cartridge” reel) is part of a removable cartridge. Upon insertion of the cartridge into the tape drive, the storage tape on the cartridge reel must be coupled to the take-up reel of the tape drive (e.g., via respective leaders). After coupling, the tape is unwound from the cartridge reel, moved past the tape head assembly and wound onto the take-up reel via a drive motor. Next, the tape is unwound from the take-up reel, moved past the tape head assembly and wound onto the cartridge. Subsequently, the storage tape must be uncoupled from the take-up reel, prior to removing the cartridge from the tape drive. In another arrangement, both reels are part of a cassette which is inserted into a tape drive and driven by a drive motor.\nMore recently, a popular trend is towards multi head, multi-channel fixed head structures with narrowed recording gaps and data track widths so that many linear data tracks may be achieved on a tape medium of a predetermined width, such as one-half inch width tape. To increase the storage density and reduce access time of magnetic tapes, data tracks on the tape are arranged with greater density and the tape is streamed by a tape head at increasingly faster rates.\nHowever, increased storage density and linear speed can lead to higher error rates when reading and/or writing on the tape due to both “axial” and “radial” run-out. Axial run-out refers to lateral or transverse motion of the tape on the reel assembly relative to a head assembly in a tape drive as the tape streams by the head assembly, and is generally defined as the peak-to-peak distance of the undesirable movement (in-plane) of the tape perpendicular to its prescribed longitudinal direction of motion past the head assembly. Radial run-out refers to tension variations in the tape in a direction perpendicular to the axis of rotation of the reel assembly.\nOften, axial and radial run-out are caused by various fabrication and assembly tolerances in the engagement between a reel assembly and a corresponding drive assembly (e.g., reel driver, drive shaft). If the engagement is imprecise because of an offset (e.g., in the axial and/or radial directions), the tape path may vary resulting in excess lateral tape motion and/or tension variations in the tape. Both axial and radial run-out can cause errors in the reading and/or writing process by limiting the degree to which a head assembly can locate a particular data track. As a result, axial and radial run-out and the ability to compensate for the same are major limiting factors in determining the minimum width of a track and the minimum spacing between tracks on the tape. As run-out is reduced, tracks may be stored more densely on the storage tape and storage tape capacity can be increased.\nThe common approach to limiting axial and radial run-out is by way of improving the fabrication and/or assembly tolerances of and/or between the reel and motor assemblies to allow for a more precise engagement between the same. In some arrangements, the reel and drive assemblies have corresponding apertures through which fasteners can be inserted and/or threaded to limit relative movement. In other arrangements, corresponding interlock components (e.g., teeth, splines, and the like) may be disposed on the reel and drive assemblies and which may respectively interlock to limit relative movement between the reel and drive assemblies. However, even relatively small mismatches in the matched features or teeth can cause run-out control and axial location difficulties. Furthermore, less than full engagement between the reel and drive assemblies can cause various types of tape loading or tension errors to occur and possibly lead to increased lateral tape motion. In this regard, the tolerance “stack-up” between the reel assembly, the drive assembly, and the like can cause various types of errors in reading and/or writing processes and hinder the pursuit of increased magnetic tape data density."}
{"text": "Network-based media services, such as internet television services and other streaming media services, may store a user profile for each user of the media service. Various settings related to media consumption and user privacy may be stored in a user's profile and adjusted by accessing the profile. For example, a user may be able to choose whether to share information related to any media content consumed by the user with other users in a social network. Such privacy settings generally apply to all media content consumed while the settings are active, and may be adjusted by accessing the user profile via a menu system.\nUser profiles also may be affected by activities carried out by a user while the user is consuming media (e.g., watching television, listening to music). For example, a media consumption history stored in a user's profile may affect recommendations of other media content items generated by the media server for that user. By using the user's activities to personalize the user profile, a user may avoid having to manually provide these user profile settings."}
{"text": "The invention relates generally to shrouds for motor vehicle cooling fans and more particularly to a shroud for a motor vehicle engine cooling fan which utilizes a Coanda surface to provide air flow through the annulus between the fan blade tips and the shroud and includes an integral air supply.\nAs motor vehicle engine compartment designs continue to evolve in response to increasing demands of vehicle and engine efficiency, operating temperatures continue to increase while the engine compartment's frontal area and natural air flow continue to reduce. All of these considerations conspire to increase underhood operating temperatures.\nNonetheless, a vehicle traveling at highway speeds at elevated ambient temperatures presents no significant engine cooling problems. Likewise, a vehicle stopped in traffic in moderate ambient temperatures presents no significant cooling problems. The combination, however, of high ambient temperatures and operation in congested, slow moving traffic wherein air heated by one vehicle is ingested by an adjacent vehicle and heated further represents an acknowledged severe engine operating condition. A second severe operating condition known as \"hot soak\" occurs when the engine has been subjected to heavy load by, for example, pulling a trailer uphill and the vehicle then stops. Operation under these conditions demands operation of and dependence upon the engine driven cooling fan. Operation in these conditions also demands the highest possible efficiency from the fan in order to achieve maximum cooling and safe engine operating conditions.\nSuch fan efficiency is achieved by well-known and recognized parameters such as the number of fan blades and their configuration as well as a properly designed radiator/fan shroud which maximizes radiator air flow and heat transfer while minimizing leakage and back flow around the fan.\nIn this regard, a problem inherent in motor vehicle design typically interferes with the attainment of high fan efficiencies. This problem results from the mounting of the radiator and fan shroud to the vehicle body whereas the fan is mounted upon the engine which is, in turn, secured to the vehicle body or frame through a plurality of engine mounts. These engine mounts are typically resilient and allow controlled motion of the engine and associated drive train components relative to the body or frame in response to engine reaction torque and vehicle acceleration and deceleration. While the spacing of the fan tips from the shroud can vary depending upon the fan and shroud location relative to the engine mount, the stiffness of the engine mounts and other variables, it has been found that spacing on the order of one-half inch (12.7 mm) to one inch (25.4 mm) or more is necessary to ensure that given the greatest excursion of the engine and fan relative to the shroud and vehicle body, the fan does not contact the shroud.\nUnfortunately, the introduction of an annular space of this size has a significant deleterious effect on fan efficiency. Fan efficiencies in such configurations have been determined to be on the order of sixteen percent. Viewed not only from the perspective of fan efficiency but also from the perspectives of achieving necessary engine cooling with a given fan size and overall engine efficiency and fuel consumption, this is not a desirable figure. Accordingly, it is apparent that improvements in the configuration of motor vehicle cooling fans which provide improved fan efficiency and thus motor vehicle cooling are desirable."}
{"text": "The present invention relates to an external cavity laser on silicon, and more specifically, to temperature insensitive external cavity lasers on silicon.\nAn optical cavity or optical resonator is an arrangement of mirrors that forms a standing wave cavity resonator for light waves. Optical cavities are a major component of lasers, surrounding the gain medium and providing feedback of the laser light. They are also used in optical parametric oscillators and some interferometers. Light confined in the cavity reflects multiple times, producing standing waves for certain resonance frequencies. The standing wave patterns produced are called modes. Longitudinal modes differ only in frequency while transverse modes differ for different frequencies and have different intensity patterns across the cross section of the beam.\nDifferent resonator types are distinguished by the focal lengths of the two mirrors and the distance between them. Flat mirrors are not often used because of the difficulty of aligning them to the needed precision. The geometry (resonator type) must be chosen so that the beam remains stable, which means that the size of the beam does not continually grow with multiple reflections. Resonator types are also designed to meet other criteria such as minimum beam waist or having no focal point inside the cavity. Optical cavities are designed to have a large Q factor, which means that the light beam will reflect a very large number of times with little attenuation. Therefore, the frequency line width of the beam is very small compared to the frequency of the laser.\nLight confined in a resonator will reflect multiple times from the mirrors, and due to the effects of interference, only certain patterns and frequencies of radiation will be sustained by the resonator, with the others being suppressed by destructive interference. In general, radiation patterns which are reproduced on every round-trip of the light through the resonator are the most stable, and these are the eigenmodes, known as the modes, of the resonator.\nResonator modes can be divided into two types: longitudinal modes, which differ in frequency from each other; and transverse modes, which may differ in both frequency and the intensity pattern of the light. The basic or fundamental transverse mode of a resonator is a Gaussian beam.\nThe most common types of optical cavities consist of two facing plane (flat) or spherical mirrors. The simplest of these is the plane-parallel or Fabry-Pérot cavity, consisting of two opposing flat mirrors. Plane-parallel resonators are therefore commonly used in microchip lasers, microcavity lasers, and semiconductor lasers. In these cases, rather than using separate mirrors, a reflective optical coating may be directly applied to the laser medium itself."}
{"text": "A number of non-inflammatory pathological conditions which can afflict warm-blooded animals involve abnormal calcium and phosphate metabolism. Such conditions may be divided into two broad categories.\n1. Conditions which are characterized by anomalous mobilization of calcium and phosphate leading to general or specific bone loss or excessively high calcium and phosphate levels in the fluids of the body. Such conditions are sometimes referred to herein as pathological hard tissue demineralizations.\n2. Conditions which cause or result from deposition of calcium and phosphate anomalously in the body. These conditions are sometimes referred to herein as pathological calcifications. Many of them result in tissue inflammation, but some do not.\nThe first category includes osteoporosis, a condition in which bone hard tissue is lost disproportionately to the development of new hard tissue. Marrow and bone spaces become larger, fibrous binding decreases, and compact bone becomes fragile. Osteoporosis can be subclassified as post-menopausal, senile, drug induced (e.g., adrenocorticoid, as can occur in steroid therapy), disease induced (e.g., arthritic and tumor), etc., however, the manifestations are essentially the same. Another condition in the first category is Paget's disease (osteitis deformans). In this disease, dissolution of normal bone occurs which is then haphazardly replaced by soft, poorly mineralized tissue such that the bone becomes deformed from pressures of weight bearing, particularly in the tibia and femur. Hyperparathyroidism, hypercalcemia of malignancy, and osteolytic bone metastases are conditions also included in the first category.\nThe second category, involving conditions manifested by anomalous calcium and phosphate deposition, includes non-inflammatory afflictions such as myositis ossificans progressiva and calcinosis universalis.\nPolyphosphonic acids and their pharmaceutically-acceptable salts have been proposed for use in the treatment and prophylaxis of such conditions. In particular diphosphonates, like ethane-1-hydroxy-1,1-diphosphonic acid (EHDP), propane-3-amino-1-hydroxy-1,1-diphosphonic acid (APD), and dichloromethane diphosphonic acid (Cl.sub.2 MDP) have been the subject of considerable research efforts in this area. Paget's disease and heterotopic ossification are currently successfully treated with EHDP. The diphosphonates tend to inhibit the resorption of bone tissue, which is beneficial to patients suffering from excessive bone loss. However, EHDP, APD and many other prior art diphosphonates have the propensity of inhibiting bone mineralization when administered at high dosage levels.\nIt is believed that mineralization inhibition is predominantly a mass related physico-chemical effect, whereas resorption inhibition results from a biological interaction with the cells. It is therefore desirable to develop more biologically potent diphosphate compounds that can be administered at low dosage levels which cause little or no mineralization inhibition, thereby resulting in a wider margin of safety. Low dosage levels are also desirable to avoid the gastro-intestinal discomfort (like diarrhea) sometimes associated with oral administration of large quantities of diphosphonates."}
{"text": "1. Field of the Invention\nThe invention generally relates to the separation of particulate materials including minerals, ores, industrial wastes and the like. More particularly, the invention deals with the separation of particulate and granular materials by utilizing a divergent electric field which effects the separation in accordance with the values of the dielectric constants of the materials being separated relative to the dielectric constant of a liquid dielectric medium within which the separation occurs.\n2. Description of the Prior Art\nThe prior art recognizes that separation of various particulate materials in dry form or solids from liquids may be achieved by utilizing electrophoresis. Electrophoresis arises from the electrostatic attraction of charged electrodes for the charged particles. The direction of motion of the charged particles is dependent on the direction of the electric field, such that reversal of the field reverses the direction of travel. Electrostatic separation is generally conducted with air as the fluid dielectric medium within which two electrodes of opposite polarity are disposed. The high voltage electric field required to effect this type of separation is not necessarily divergent and, in most cases, is relatively uniform. It is generally accepted that optimum electrostatic separations are effected when the dielectric constant of the electrode gap be as near unity as possible.\nThe Tyson U.S. Pat. No. 2,485,335 is an example of a typical prior art system for separating materials by means of an electrostatic field. Tyson teaches that solid polymer particles may be separated from an associated reaction liquid by passing the slurry of liquid and particles through an electrostatic field created by an electrical discharge potential. The electrical discharge causes the polymer particles to adhere to the surface of the rotating drum which carries a charge opposite to that carried by the particles. The liquid is caused to drain away and the polymer particles adhering to the drum are subsequently removed therefrom.\nBy contrast, separation of materials through dielectrophoresis occurs because of the tendency of matter to become polarized in a non-uniform electric field and move into regions of highest field strength or intensity. Dielectrophoresis only occurs in a divergent or non-uniform electric field created in a fluid dielectric medium. In a divergent electric field, particles in a less-polar suspension migrate to and accumulate where the field is strongest. This movement does not require the particles to be charged but depends instead on the force experienced by all polarizable materials in non-uniform electric fields. The uncharged particle is believed to be polarized by the applied voltage and, consequently, moves due to the unequal pulls exerted by the forward and backward direction of the non-uniform electric field. The most polar material moves toward the area of greatest field intensity, independent of the direction of the field and the type of applied voltage, whether such voltage is in the form of alternating current, direct current or pulsating direct current. It is therefore recognized that for separation to occur, the dielectric constant of the particles to be separated must be higher than that of the liquid dielectric medium within which the divergent electric field is created.\nThe Walker U.S. Pat. No. 3,304,251 exemplifies a prior art system for separating particles by the utilization of dielectrophoresis or a non-uniform electric field. This is achieved by passing an inner wire electrode through an outer tubular electrode within which a contact zone is provided for the materials being separated. A liquid dispersion of the particles is contacted in the zone with a non-uniform electrical field such that the particles migrate and adhere to the wire electrode which removes the particles from the contact zone."}
{"text": "Phosphodiesterase type 5 inhibitors are drugs used to block the degradative action of phosphodiesterase type 5 on cyclic GMP in the smooth muscle cells lining the blood vessels supplying the corpus cavernosum of the penis. These drugs are commonly used in the treatment of erectile dysfunction.\nPhosphodiesterase type 5 inhibitors are commonly delivered orally. Currently, no transdermal formulations of phosphodiesterase type 5 inhibitors have been approved by the FDA. Accordingly, systems and methods for transdermally delivering clinically useful amounts of phosphodiesterase type 5 inhibitors are needed."}
{"text": "1. Field of the Invention\nThis invention relates to downhole drilling, and more particularly to apparatus and methods for transmitting data along a downhole drill string.\n2. Description of the Related Art\nFor half a century, the oil and gas industry has sought to develop downhole telemetry systems that enable high-definition formation evaluation and borehole navigation while drilling in real time. The ability to transmit large amounts of sub-surface data to the surface has the potential to significantly decrease drilling costs by enabling operators to more accurately direct the drill string to hydrocarbon deposits. Such information may also improve safety and reduce the environmental impacts of drilling. This technology may also be desirable to take advantage of numerous advances in the design of tools and techniques for oil and gas exploration, and may be used to provide real-time access to data such as temperature, pressure, inclination, salinity, and the like, while drilling.\nIn order to transmit data at high speeds along a drill string, various approaches have been attempted or suggested. One approach that is currently showing promise is to incorporate a “network” of data transmission cable and other communication equipment into the drill string. Due to the length of drill strings, which may exceed 20,000 feet, such a network may require placing network “nodes” at selected intervals along the drill string. These nodes may act as repeaters to amplify the data signal and provide points of data collection along the drill string. Communication elements, such as magnetic couplers, may be incorporated into the ends of downhole tools to transmit data across the tool joints. Transmission lines, such as electrical cables, may be incorporated into the downhole tools to transmit data therealong.\nUnfortunately, unlike conventional above-ground networks, a downhole network is constrained by the physical limitations of the downhole drill string. In particular, a downhole drill string is a linear structure, making it very difficult to build redundancy (and thereby reliability) into the downhole network. As a result, any break or malfunction in the data transmission path along the drill string may cause communication to be lost between the surface and downhole components. Because the drill string may include many hundreds of downhole components (e.g., sections of drill pipe, drill collar, bottom-hole assembly components, etc.), a single break or malfunction in any downhole component can break the communication path and cause the network to lose much if not all of its functionality.\nIn view of the foregoing, what are needed are apparatus and methods to provide multiple redundant paths of communication in a downhole network. Such apparatus and methods may be used to significantly improve the reliability of downhole communication networks."}
{"text": "With a technology for adaptive speech speed conversion, for a given playback speed, such as 1× speed (playback in real time) or 2× speed (playback in half of real time), the speed is not changed by a uniform factor α over the entire input signal, but rather the speed is changed in each section by a factor larger or smaller than the factor α so as to balance the overall playback time to be the same as when speech speed is converted at the uniform factor α. Thereby it is aimed to generate speech speed converted voice that is “slower and easier to hear” for the listener than when speech speed is converted at the uniform factor α.\nSome techniques for achieving the above include (1) lowering the speech speed where the fundamental frequency is high and raising the speech speed where the fundamental frequency is low, (2) treating an interval spoken in one breath as a unit, lowering the speech speed at the start of the interval, and gradually raising the speech speed towards the end of the interval in accordance with changes in the fundamental frequency, and (3) shortening a silent interval between intervals spoken in one breath to a degree that preserves a natural sound (for example, see Patent Literature 1).\nAnother technique treats a silent interval of at least a given length as a pause, and in a voice interval located between pauses, lowers the speech speed at the start of the voice interval, progressively raises the speech speed during a given time T based on a predetermined decreasing function, and after the given time T elapses, changes the factor for lowering the speech speed by taking into consideration the relative magnitude of the maximum fundamental frequency in each voice interval (for example, see Patent Literature 2).\nWithin the speech speed control disclosed in Patent Literature 1 or Patent Literature 2, another known technique allows for a brief silent interval within a voice interval located between pauses to be shortened to a degree that still preserves a natural sound. This technique also lowers the subsequent speech speed in so far as possible when the speech speed of each section matches, or is only slightly later than, the time assumed when the speech speed being converted at a uniform factor α, and reduces the amount by which the subsequent speech speed is lowered as the speech speed of each section is increasingly later than the time assumed when the speech speed is converted at the uniform factor α. This technique thereby lessens misalignment, in so far as possible, with the time assumed when the speech speed of each section of the speech speed converted voice is converted at the uniform factor α (for example, see Patent Literature 3).\nFurthermore, when separating the input signal into voice intervals and silent intervals, lowering the speech speed in the voice intervals, and shortening the silent intervals, the output voice length extends beyond the input signal length per unit time due to the lowering of the speech speed in the voice intervals. It thus becomes necessary to store the voice after speech speed conversion temporarily in memory, yet there is a limit on memory capacity. Therefore, a technique is known for gradually raising the speech speed in the voice intervals and increasing the amount cut from the silent intervals in accordance with the remaining memory capacity (for example, see Patent Literature 4 and 5).\nAdditionally, a technique is known for determining the speech speed of each section using a coefficient such that the speech speed is inversely proportional to the increase or decrease of the magnitude (power) or pitch (fundamental frequency) of the input signal, or a coefficient such that the speech speed is inversely proportional to the nth power of the value of the magnitude or volume of the input signal (for example, see Patent Literature 6)."}
{"text": "One problem associated with known injection molds is that usually a two part mold is configured in such a way that a product formed in the mold is not held in either part of the two part mold so that the product must be fixed to a machine by particular holding, gripping or fixation devices for further machining. Thus, there remains a need for an improved method for manufacturing a product by injection molding with an injection mold which permits machining or a removal of a portion of the injection mold so that the product is firmly held in the remaining part of the injection mold and can be further processed or machined without the need for additional holding, gripping or fixation devices."}
{"text": "German Published Patent Application No. 197 14 784 describes a compact drive, which includes an electric motor, at whose one end face a gear unit is arranged, and at whose other end face a frequency converter is arranged. The electronics region and the motor region must be sealed with respect to the gear unit. In this context, it is disadvantageous that the axial length is large and that a power take-off can only be provided at one end face of the compact drive.\nJapanese Published Patent Application No. 2002-336305 describes a gear motor, which does not include a frequency converter, however. Between the gear unit, motor, and brake, as well as between the corresponding housing parts, there are also interfaces that must be precisely manufactured and are therefore complex and expensive.\nA frequency converter, which is axially mounted to an electric motor behind its blower and is therefore cooled by its air stream, is described in German Published Patent Application No. 196 22 396. However, this therefore requires a large amount of space and, in addition, a fan that must be able to direct ambient air past.\nA piezoelectric brake is described in European Published Patent Application No. 0 694 203.\nGerman Published Patent Application No. 102 07 760 describes an adjusting gear, which, however, requires a large amount of space and an interface to the rotor of the motor.\nGerman Published Patent Application No. 198 48 324 describes a coupling, which, however, requires additional, complicated and expensive machining of the rotationally mounted parts. In particular, not only are the sun and the rotor shaft included, but also a coupling sleeve. A higher number of parts also means higher storage costs.\nEuropean Published Patent Application No. 1 081 827 describes an electric tool, in which a gear unit can be driven by an electric motor. However, an interface, which is difficult to manufacture, is provided between the housing parts of the motor and the gear unit.\nAn axially offset right-angle drive, i.e., a spiroid gear unit, which is connectable to an electric motor and can be driven by it, is described in German Published Patent Application No. 43 09 559."}
{"text": "1. Field of the Invention\nThe present invention relates to a thin-film magnetic head for reading and writing signals by means of a perpendicular magnetic recording, a head gimbal assembly (HGA) with the thin-film magnetic head and a magnetic disk drive apparatus with the HGA, and to a manufacturing method of the thin-film magnetic head for the perpendicular magnetic recording.\n2. Description of the Related Art\nIn recent years, the perpendicular magnetic recording has been actively developed instead of conventional longitudinal magnetic recording to realize more improvement in areal recording density of the magnetic disk drive apparatus.\nIn the perpendicular magnetic recording, a demagnetization field drastically decreases in a magnetization transition region, and therefore, the magnetization transition width can become much smaller than that of the longitudinal magnetic recording. Furthermore, the perpendicularly recorded magnetization is not greatly affected by a thermal fluctuation that becomes serious problem in the longitudinal magnetic recording with higher recording density. Therefore, it becomes possible to obtain stable higher-density recording in the perpendicular magnetic recording.\nA thin-film magnetic head for the perpendicular magnetic recording has been proposed, which has a single pole structure including a main magnetic pole, an auxiliary magnetic pole as a return yoke, and an inductive coil inducing a magnetic flux in these magnetic poles. Generally, the magnetic head also has magnetic shields provided in the leading and/or trailing sides of the main magnetic pole. As the same time, a magnetic disk has been proposed which has a multilayer structure of a soft-magnetic backing layer acting as a part of a magnetic circuit and a perpendicular magnetic recording layer.\nIn the thin-film magnetic head, a field generated from the main magnetic pole performs writing to the magnetic disk. FIG. 1 shows a schematic drawing illustrating a simulated distribution of a magnetic field generated from the conventional main magnetic pole, and recorded bits. In the figure, only the magnetic field generated from the main magnetic pole 10 is illustrated. The cross-section of the main magnetic pole 10 has a trapezoid shape due to the formation of bevel angles against a skew angle arising from an arm driving by means of a rotary actuator, which possibly causes unnecessary signals to be recorded on adjacent tracks. Broken lines show magnetic field contour lines.\nAs shown in FIG. 1, the contours of the magnetic field generated from the main magnetic pole 10 has a curvature in such a way as to surround the main magnetic pole 10 itself. Recorded bits 11 are formed on the magnetic disk along the contour line that lies in the trailing side, that is, in the downstream side of the disk rotation, and is equivalent to a coercive force of the magnetic disk. As a result, the shape of the magnetization transition region 12 as a boundary region between recorded bits 11 also has a curvature. Therefore, when a magnetoresistive (MR) read head element reads signal fields generated from the recorded bits 11 with the curved transition regions 12, a jitter becomes larger due to the increase in a reverse width of a reproduction power, and the larger jitter causes an error rate to be increased. Moreover, the length of the recorded bit 11 can not be shorten down to the curvature width WTC of the transition region 12, which hinders higher recording density.\nAs a technique improving the magnetization transition region, for example, a structure that an opposed portion of the auxiliary magnetic pole to the main magnetic pole is protruded in order to suppress the broadening of a magnetic field distribution at track edges, is described in Japanese Patent Publication No. 2002-092820A, in which no curvature of the magnetic field is mentioned. Further, in U.S. patent application Ser. No. 10/880,509, it is described that the main magnetic pole has a double-layered structure of the first layer formed of, for example, CoFe and the second layer formed of, for example, NiFe, CoNiFe or CoFe, though not to improve the magnetization transition region. Furthermore, a structure of the main magnetic pole that has a concave portion in the trailing side of an air bearing surface (ABS) in order to reduce the degree of the curvature of the transition region, is described in U.S. Pat. No. 6,741,421.\nHowever, the structure with the protrusion described in Japanese Patent Publication No. 2002-092820A is not adequate to control the degree of the curvature. The structure is just for suppressing the broadening of the magnetic field distribution and concentrating the magnetic field, and can not positively control the curvature. Further, the double-layered structure described in U.S. patent application Ser. No. 10/880,509 has no purpose of improving the transition region. Actually, there is no specification about the thickness of each layer in the description, and mere double-layered structure itself can not realize the desired degree of the curvature.\nFurthermore, the structure with the concave portion described in U.S. Pat. No. 6,741,421 has a difficulty in forming the concave portion with high accuracy. In other words, the size and depth of the concave portion can not be precisely controlled because the concave portion is formed by digging the magnetic material layer with ion-milling method and so on. However, the degree of the curvature of the magnetic contour, which determines the shape of the magnetization transition region, is sensitive to the shape of the concave portion. Therefore, the contour with small curvature can not be obtained stably.\nFurthermore, because a non-magnetic material is embedded in the concave portion, an adequate write field is difficult to be obtained near the trailing edge of the main magnetic pole, which is a main portion for writing, compared to the case without the concave portion.\nMoreover, in an MR read head element utilizing giant magnetoresistive (GMR) effect or tunnel magnetoresistive (TMR) effect, a pinned layer, a separation layer and a free layer, which are main components, are usually stacked sequentially on a hard-bias layer. Therefore, in some cases, the free layer that receives signal fields has a curvature with a reverse direction to that of the curvature of the transition region. On this case, a magnetic sensitivity contour that expresses a sensitivity level to a magnetic field of the MR read head element also has a curvature with the reverse direction to that of the curvature of the transition region. As a result, a jitter becomes larger due to the increase in the reverse width of the reproduction power, and the larger jitter causes the error rate to be increased.\nTherefore, only reducing the degree of the curvature of the magnetic field contour is not sufficient, and it is needed to design the magnetic field contour with wide design range including the direction of the curvature in consideration of the magnetic sensitivity contour of the MR read head element. Without such a design, the error rate can not be sufficiently decreased by reducing the reverse width of the reproduction power."}
{"text": "Field of the Invention\nThe present invention relates to a magnetic recording medium, a magnetic signal reproduction device and a method of manufacturing a magnetic recording medium.\nDiscussion of the Background\nIn a particulate magnetic recording medium, a magnetic layer containing ferromagnetic powder and binder is present on a nonmagnetic support and a nonmagnetic layer containing nonmagnetic powder and binder is optionally present between the nonmagnetic support and the magnetic layer.\nConventionally, nonmagnetic powder (nonmagnetic particles) that can contribute to reducing the coefficient of friction during running has been incorporated along with ferromagnetic powder into the magnetic layer of a particulate magnetic recording medium (for example, see Japanese Unexamined Patent Publication (KOKAI) No. 2011-48878 or English language family members US2011/052908A1 and U.S. Pat. No. 8,535,817, which are expressly incorporated herein by reference in their entirety."}
{"text": "A plated material produced by providing a plating layer of, for example, tin (Sn) or a tin alloy, on an electroconductive base material, such as copper (Cu) or a copper alloy (hereinafter, appropriately referred to as base material), is known as a high performance conductor material having the excellent electroconductivity and mechanical strength of the base material, as well as the excellent electrical connectivity, corrosion resistance, and solderability of the plating layer. Thus, such plated materials are widely used in various terminals, connectors, and the like.\nIn recent years, since a fitting-type connector is multipolarized with advancement of electronic control, a considerable force is necessary for plugging a group of male terminals into/out of a group of female terminals. In particular, plugging-in/out such a connector is difficult in a narrow space such as the engine room of a vehicle, and it has been strongly demanded to reduce the force for plugging in/out such a connector.\nIn order to reduce the plugging-in/out force, the Sn plating layer on the surface of the connector terminal may be thinned to weaken contact pressure between the terminals. However, because the Sn plating layer is soft, a fretting phenomenon may occur between contact faces of the terminals, thereby causing inferior conduction between the terminals.\nIn the fretting phenomenon, the soft Sn plating layer on the surface of the terminal wears and is oxidized, becoming abrasion powder having large specific resistance, due to fine vibration between the contact faces of the terminals caused by vibration and changes in temperature. When this phenomenon occurs between the terminals, conduction between the terminals results in inferior. The lower the contact pressure between the terminals, the more the fretting phenomenon is apt to occur.\nPatent Literature 1 describes an electrically conductive material for a connecting part, having a Cu—Sn alloy coating layer and a Sn coating layer, formed in this order, on the surface of a base material formed from a Cu strip, wherein the Cu—Sn alloy coating layer has the exposure area ratio at the material surface of 3 to 75%, the average thickness of 0.1 to 3.0 μm, and the Cu content of 20 to 70 at %; and the Sn or Sn alloy coating layer has the average thickness of 0.2 to 5.0 μm. It is also described that a Cu—Sn alloy coating layer is formed by performing a reflow treatment.\nAccording to Patent Literature 1, when this electrically conductive material is used in, for example, a multipole connector in automobiles, a low insertion force upon fitting of male and female terminals is attained, and the assembly operation can be efficiently carried out; and the electrically conductive material is considered to be able to maintain electrical reliability (low contact resistance), even if maintained for a long period of time under a high temperature atmosphere, or even under a corrosive environment.    Patent Literature 1: JP-A-2006-77307 (“JP-A” means unexamined published Japanese patent application)"}
{"text": "The present invention relates to fins and their attachment to small water craft and, more particularly to a system for attaching a variety of fin configurations to a surfboard.\nAs with many types of sporting equipment, many variations on the basic article have been developed to tailor the article to the preferences of the user and the conditions of use. So it is with the fins attached to the rear undersides of surfboards.\nIt is an obvious advantage for the owner of a board to be able to interchange the fin or fins both to experiment, with the aim of finding the fin most suited to their body weight and style of riding, and from time to time to suit the conditions in which the board is to be ridden.\nInterchangeable fin systems are known and comprise of a selection of fins retained in a so-called fin box. Particular examples may be found in U.S. Pat. No. 5,830,025, U.S. Pat. No. 5,975,974 and WO 01/70565. The fin box is generally a permanent fixture built into the body of the board during its manufacture and, again generally includes a slot or slots into which the tab or heel of the fin may be inserted. The means for retaining an inserted fin securely in the fin box are many and varied, ranging from relatively simple snap-in systems to arrangements incorporating fasteners requiring the application of tools to install and remove a fin from its fin box.\nA disadvantage with these systems is that the fin box is designed to take only one configuration of a fin tab, so that the owner of a board who wishes to change or experiment with a variety of fins from manufacturers other than the manufacturer of his or her board, or at least the fin boxes fitted to the board, is restricted in choice.\nIt is an object of the present invention to address or at least ameliorate some of the above disadvantages."}
{"text": "Field of the Invention\nThe present invention concerns the reconstruction of a magnetic resonance (MR) image from radial collection of MR data in a manner that avoids artifacts (in particular, artifacts caused by metal parts).\nDescription of the Prior Art\nIn magnetic resonance imaging, raw data (magnetic resonance measurement data) are acquired from a subject by the operation of a magnetic resonance scanner in which the subject is situated. The acquired raw data are entered as complex numbers at respective data points of a memory. This collection of raw data in the memory is known as k-space data. The entry of the raw data at the respective data points of the memory is known as “scanning” k-space, or scanning a k-space memory. The path composed of data points along which the raw data are successively entered in k-space is known as a “k-space trajectory.”\nThe k-space data are then transformed into image data, typically by operation of a Fourier transform on the k-space data. The k-space trajectory that is used to enter the raw data into k-space may have an influence on the type or degree of artifacts that may occur in the image represented by the image data. Trajectories such as a Cartesian trajectory, which proceeds along a rectilinear path, and a radial trajectory, which proceeds radially outwardly from a center of k-space to a periphery of k-space, are among the known types of k-space trajectories.\nThe radial scanning of k-space, which is also referred to as the radial collection of MR data, is a known method, for obtaining T1-weighted MR data, for example, which is used for musculoskeletal imaging (imaging of the musculoskeletal system.\nCompared with the traditional Cartesian collection of MR data, the radial collection of MR data has the advantage of being less sensitive to movement. However, the radial collection of MR data is more sensitive to system failures such as gradient delays and off-resonance effects. For example, sensitivity with regard to artifacts caused by metal parts is more pronounced in the case of the radial collection of MR data than in the case of the Cartesian collection of MR data."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus to control the ratio of air to fuel of the air-fuel mixture which is supplied to an internal combustion engine, and more particularly to an apparatus provided for controlling the amount of air supplied to the main system fuel passage and the slow system fuel passage of a carburetor for supplying an air-fuel mixture to an internal combustion engine by means of a single proportional control solenoid valve.\n2. Description of the Prior Art\nIt is known to incorporate a catalytic converter have ternary catalytically active substances into the exhaust system of the internal combustion engine of an automobile in order to simultaneously reduce the injurious components of the exhaust gas, such as hydrocarbons (HC), carbon monoxide (CO) and nitrogen oxides (NOx). It is required to supply an air-fuel mixture of an air-fuel ratio corresponding to the stoichiometric air-fuel ratio into the cylinders of an internal combustion engine, since the maximum cleaning efficiency of the ternary catalytically active substances is attained when the exhaust gas is produced by the combustion of an air-fuel mixture of the stoichiometric air-fuel ratio. An air-fuel ratio controller to meet such a requirement is proposed, which includes a main system air-bleeding passage connected to the main system fuel passage connecting to the main jet of a carburetor, a slow system air-bleeding passage connected to the slow system fuel passage connecting to the slow system fuel supply port of the carburetor, and a main system proportional control solenoid valve for controlling the main air bleeder and a slow system proportional control solenoid valve for controlling the slow air bleeder, which are disposed within the main system air-bleeding passage and the slow system air-bleeding passage, respectively, and are adapted to be controlled by a control signal provided by converting the output signal of an oxygen sensor disposed within the exhaust passage of the engine by means of an electronic control unit. It is proposed to control the rate of air to be supplied to the main system fuel passage and to the slow system fuel passage through the main system air-bleeding passage and through the slow system air-bleeding passage, respectively, with an air-fuel ratio controller as described above, so that the air-fuel ratio of the air-fuel mixture which is supplied into the cylinders of an engine is controlled so as to be close to the stoichiometric air-fuel ratio.\nHowever, heretofore known air-fuel ratio controllers of this type for carburetor are provided with individual proportional control solenoid valves in the main system air-bleeding passage and the slow system air-bleeding passage, respectively, and are adapted to control both proportional control solenoid valves simultaneously with the control signal of a single system provided by the control signal producing circuit of an electronic control unit. Accordingly, the electronic control unit is required to supply an electric current simultaneously to the main system and the slow system proportional control solenoid valves. Therefore, the electronic control unit is necessary to be capable of supplying an electric current twice as much as that to be supplied to a single solenoid. Thus such conventional air-fuel ratio controllers have a disadvantage that the control signal producing circuit of the electronic control unit must consist of elements which have superior current capacities."}
{"text": "The present invention relates to a seat sliding device for sliding a seat of a vehicle.\nA kind of seat sliding device for sliding a seat of a vehicle as shown in FIGS. 1 and 2 has been disclosed. A pair of these kind of seat sliding devices are provided for a seat and arranged in parallel with each other in a back-and-forward direction of the seat. The seat sliding device has a pair of rails, a lower rail and an upper rail. The lower rail is fixedly mounted on a vehicle body and the upper rail is fixedly mounted on the seat and slidably engaged with the lower rail by rolling members.\nA horizontally settled seat sliding device is shown in FIG. 1, and a longitudinally settled seat sliding device is shown in FIG. 2.\nIn the horizontally settled seat sliding device, a lower rail 1 mounted on a vehicle body is provided with a lower base 1a which has a somewhat U shaped cross section and a pair of lower flanges 1b which are provided on both upper ends of the lower base 1a and extend toward the outside thereof. The upper rail 2 is directly mounted on a bottom face of the seat, or it can be indirectly mounted on the bottom face of the seat through a supporting member. Further, the upper rail 2 comprises an upper body 2a (which covers an open end of the lower rail 1) and a pair of upper flanges 2b. These upper flanges 2b are provided on both side ends of the upper body 2a, extend downwards, and are engaged with the pair of lower flanges 1b, as shown in FIG. 1. The upper rail 2 is supported by a roller 3 which is provided inside the lower rail 1. A pair of balls 4 are provided in a pair of spaces which are formed by the lower flanges lb and the engaged upper flanges 2b. These balls receive a pulling force applied to the upper rail 2 in order to slide the upper rail 2 smoothly along the lower rail 1. Further, the balls 4 and roller 3 are retained by a retainer so that all these rolling members are moved together inside a seat sliding body M which is formed by the upper rail 2 and lower rail 1.\nFIG. 2 shows a longitudinally settled seat sliding device in which a lower rail 6 and an upper rail 7 are longitudinally settled. In direct contrast to the above described horizontally settled seat sliding device, the lower rail 6 has a structure identical to that of the upper rail 2, and the upper rail 7 has a structure identical to that of the lower rail 1. The other members such as the roller 3, the balls 4, and the retainer 5 in FIG. 2 correspond to those shown in FIG. 1 which have the same reference numbers. An N depicts a seat sliding body of the longitudinally settled seat sliding device.\nIn the above-described horizontally settled seat sliding device, since the cross section of the upper rail 2 has a somewhat rectangular shape as shown in FIG. 1 and the longitudinal length of the upper rail 2 is much shorter than the horizontal length thereof, a longitudinal flexural strength of the upper rail 2 is insufficient. Therefore, when the plate thickness of the upper rail 2 is thickened in order to strengthen the flexural strength thereof, the upper rail 2 becomes heavy and results in a high cost.\nIn the case of the longitudinally settled seat sliding device, the longitudinal length of the upper rail 7 is wide enough to create a sufficient longitudinal flexural strength in the upper rail 2. However, since the horizontal width of the upper rail 7 cannot be made large enough and the upper rail 7 does not have any support means which work to support the upper rail 7 in a horizontal direction, the horizontal strength of the upper rail 7 may not be sufficient thereby leading to a possible rattly condition in which the upper rail 7 cannot slide smoothly along the lower rail 6.\nMoreover, since the structure of the lower rail 1 is not similar to that of the upper rail 2 and the structure of the lower rail 6 is not similar to that of upper rail 7, errors when forming the lower rail 1 or 6 and the upper rail 2 or 7 may occur resulting in a misfitting of the lower rail 1 in the upper rail 2 or of the upper rail 7 in the lower rail 6. In this situation, the upper rail 2 or 7 cannot slide smoothly along the lower rail 1 or 6, respectively."}
{"text": "Superconducting wires are widely used for instance in NMR (Nuclear Magnetic Resonance) devices and MRI (Magnetic Resonance Imaging) devices. Rare earth-based high-temperature superconducting wires (REBCO superconducting wires) have been produced in recent years, and research and development is being conducted in which such superconducting wires are utilized.\nJoining of superconducting wires to each other is an essential technical issue with a view for instance to achieving longer wires and realizing a persistent current mode. Low-resistance connection between REBCO superconducting wires is conventionally accomplished by soldering; as is known, however, electrical resistance values of about 10−8Ω are incurred in this technique.\nResearch has been carried out in recent years for directly joining REBCO superconducting wires to each other. PTL 1 proposes a technique relying on a melt diffusion method. In this technique, the protective layers of two REBCO superconducting wires are removed, and the wires are heated up to the melting point of superconducting layers, while causing the superconducting layers to abut each other and while applying pressure to the latter; thereby, part of the abutting superconducting layers in the thickness direction undergoes melt diffusion, and the superconducting wires become joined to each other as a result. PTL 1 reports that the REBCO superconducting wires can be connected to each other with yet lower resistance than is the case when the wires are joined by soldering."}
{"text": "1. Field of the Invention\nThe present invention is directed to process controllers. More specifically, the present invention is directed to a process controller for controlling a process in response to the duration of a controller disabling power outage."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of displaying technology, and in particular to a method for manufacturing a TFT (Thin-Film Transistor) backplane and a structure of a TFT backplane.\n2. The Related Arts\nFlat panel displays have a variety of advantages, such as thin device body, low power consumption, and being free of radiation, and are thus widely used. Currently available flat panel displays generally include liquid crystal displays (LCDs) and organic light emitting displays (OLEDs)\nOrganic light emitting displays (OLEDs) have a variety of superior properties, such as being self-luminous, requiring no backlighting, high contrast, reduced thickness, wide view angle, fast response, applicability to flexible panels, wide range of operation temperature, and having simple structure and manufacturing process, and are regarded as emerging technology of the next generation flat panel displays.\nThe TFTs are a vital constituent part of the flat panel displays. Since the TFTs can be formed on a glass substrate or a plastic substrate, they are commonly used as switch devices and driver devices for devices, such as LCDs, OLEDs, and electro-phoretic displays (EPDs).\nOxide semiconductor TFT technology is the most attention-attracting technology currently. The oxide semiconductors have a relatively high electron mobility and, compared to low-temperature poly-silicon, the oxide semiconductors have a simple manufacturing process and high compatibility to amorphous silicon manufacturing processes, being applicable to the fields of LCDs, OLEDs, and flexible displays, and are also compatible to high generation manufacturing lines, being applicable to large-, medium-, and small-sized displays, making them possess excellent prosperous future of development.\nA conventional structure of an oxide semiconductor TFT backplane that is considered relatively mature is one that includes an etch stop layer. As shown in FIGS. 1-10, a method for manufacturing a conventional oxide semiconductor TFT backplane comprises the following steps:\nStep 1: providing a substrate 100, forming a first metal layer on the substrate 100, and applying a photolithographic operation to patternize the first metal layer so as to form a gate terminal 200 on one side portion of the substrate 100 and a first metal electrode M1 on an opposite side portion of the substrate 100;\nStep 2: forming a gate insulation layer 300 on the gate terminal 200, the first metal electrode M1, and the substrate 100 and applying a photolithographic operation to patternize the gate insulation layer 300 to form a gate insulation layer via 310 for exposing a portion of the gate terminal 200;\nStep 3: forming a film on the gate insulation layer 300 and applying a photolithographic operation to patternize the film to form an island-like oxide semiconductor layer 400;\nStep 4: forming an etch stop layer 500 on the oxide semiconductor layer 400 and the gate insulation layer 300 and applying a photolithographic operation to patternize the etch stop layer 500 to a plurality of etch stop layer vias 510 for exposing portions of the oxide semiconductor layer 400;\nStep 5: forming a second metal layer on the etch stop layer 500 and applying a photolithographic operation to patternize the second metal layer to form source/drain terminals 600 on one side portion of the substrate 100 and a second metal electrode M2 on an opposite portion of the substrate 100, wherein the source/drain terminals 600 fill up the plurality of etch stop layer vias 510 to connect to the oxide semiconductor layer 400 and the source/drain terminals 600 fill up the gate insulation layer via 310 to connect to the gate terminal 200;\nwherein the first metal electrode M1, the second metal electrode M2, and a portion of the gate insulation layer 300 and a portion of the etch stop layer 500 that are sandwiched between the first and second metal electrodes M1, M2 form a storage capacitor C;\nStep 6: forming a passivation protection layer 700 on the source/drain terminals 600 and the second metal electrode M2, followed by patternizing by applying a photolithographic operation;\nStep 7: forming a planarization layer 800 on the passivation protection layer 700, followed by patternizing by applying a photolithographic operation;\nStep 8: forming a pixel electrode layer 900 on the planarization layer 800, followed by patternizing by applying a photolithographic operation;\nStep 9: forming a pixel definition layer 1000 on the pixel electrode layer 900 and the planarization layer 800, followed by patternizing by applying a photolithographic operation; and\nStep 10: forming a spacer pillar 1100 on the pixel definition layer 1000.\nThe conventional oxide semiconductor TFT backplane manufacturing method suffers certain problems, which are generally presented in three aspects: The first one is that the manufacture of the oxide semiconductor TFT backplane requires ten processes of photolithographic operation, wherein the manufacture of the etch stop layer 500 requires a complete process of photolithographic operation (including the steps of film forming, yellow light, etching, and stripping). This leads to an extended operation process, a reduced manufacturing efficiency, and an increase of manufacturing cost, and the more the manufacturing steps, the more the yield problems there will be. The second one is that the gate insulation layer 300, the oxide semiconductor layer 400, and the etch stop layer 500 are not formed consecutively so that the interfaces of the oxide semiconductor layer 400 and with respect to the other two layers may be readily contaminated by etching solutions and stripping solutions, leading to a potential risk of deterioration of the performance of the TFT. The third one is that the storage capacitor C is formed of the first metal electrode M1, the second metal electrode M2, and a portion of the gate insulation layer 300 and a portion of the etch stop layer 500 sandwiched between the first and second metal electrodes M1, M2 and due to an additional thickness resulting from the presence of the etch stop layer 500, the storage capacitor C requires an enlarged area, which causes reduction of aperture ratio.\nThus, it is desired to improve the conventional method to eliminate the problems existing therein."}
{"text": "The present disclosure relates generally to primary air inlets to supply primary air to a gas burner to facilitate combustion of fuel gas supplied to the burner. The primary air is mixed with the fuel gas prior to ignition of the gas to form a fuel-rich mixture for ignition. The remaining air required for complete combustion is obtained from the ambient air in the room following ignition, and is referred to herein as secondary air."}
{"text": "1. Field of the Invention\nThis invention relates to current motors capable of making use of stream and river water flow as well as tidal water flow for the production of power. The power so produced may be used to produce electricity from a generator or it may be used as mechanical power to convey water and the like via belt or chain mechanisms or the like.\n2. Description of the Prior Art\nThe application of hydroelectric power for the generation of electricity has been known for many years. While usually quite satisfactory, hydroelectric power installations are relatively expensive and require a long lead time to build. Furthermore, they are not always practical except in rivers of certain minimal water flow. In fact, they are not economically feasible for use in many moderate sized streams. Being fixed installations, they are virtually impossible to move and, hence, are virtually useless to meet the requirements of emergency power needs.\nIn harnessing tidal energy, as opposed to that of river flows, many suggestions have been made over the years. Many of these suggestions have entailed relatively complicated, complex networks of coffer dams, concentrating valves, storage basins and the like, all of which require relatively high expenditures and fixed installations requiring a long lead time.\nWater wheels or other overshot paddle assemblies for producing power have long been known, as well as turbine systems wherein the working fluid passes axially through rotating turbine blades.\nWith the recent renewed interest in energy production arising from the runaway costs of conventional fossil fuels, there has arisen renewed interest in rivers, streams and tides as sources of power. For example, in The Philadelphia Inquirer dated June 5, 1973, there appeared an article concerning harnessing the energy of the huge Gulf Stream off the coast of Miami, Fla. Therein it was stated that the Gulf Stream could produce as much electric power as a large nuclear generator. The solution to harnessing such power, according to that article, would be to locate large, slow-turning turbines, like giant underwater windmills beneath the surface. Calculations showed that water flows in the upper layers of the Gulf Stream could produce about 0.8 kilowatt per square meter of stream cross-section or about 25,000 megawatts. If only 4 percent of this power could be extracted, this 1000 megawatts would be equivalent to the output of a large nuclear generating station.\nMy prior U.S. Pat. No. 3,807,890, utilizes a raceway in which is disposed an inlet gate for concentrating and increasing the velocity of the water flowing through the raceway and also for directing the water to one side of the raceway for engagement with paddles of a vertically disposed water wheel. The inlet gate has a side wall which diverges against the direction of water flow. This side wall, which concentrates the water flowing into the raceway, is pivoted so that it can swing back to a neutral position to permit water flowing out of the raceway to exit freely. A second gate is disposed at the opposite end of the raceway to perform a similar function in the opposite direction through the raceway.\nU.S. Pat. No. 1,034,940 discloses a current motor having a supporting frame, cable supporting drums revolubly mounted in the frame, the drums comprising heads each of which consists of a hub having radially extending arms, a series of spokes secured at their inner ends to the arms, annular bands connecting the outer ends of the spokes and segmental plates secured to the bands to form the outer annular surface of the drums, annular channel shaped cable receiving rings secured to the annular surface of the drums, operating cables engaged with the channel shaped rings whereby the drums are operatively connected together and driven, a series of blades having a loose connection with the cables, a braking mechanism adapted to be engaged with one of the drums whereby the speed of the motor is retarded, and a power transmittng gear fixed to the coacting drum whereby the power of the motor may be transferred.\nU.S. Pat. No. 2,161,215 discloses a water current motor which included a buoyant body designed to collect and from a sluice-way for the moving water of a stream and which is capable of being moored or moved from place to place in a stream.\nAn object of the present invention is to provide an improved apparatus for utilizing the energy of flowing water, including stream, tide and river flows.\nAnother object of this invention is to provide an improved apparatus for utilizing the energy of flowing water, which apparatus is capable of being moved easily from location to location to supply, for example, energency power needs.\nStill another object is to obviate many disadvantages of prior water-energy systems."}
{"text": "Light-emitting devices that use a light-emitting diode (LED) element are widely used. Older-type light-emitting devices are configured by, as illustrated in FIG. 4, joining an LED element 33 to a substrate 31 with a die bond adhesive 32, wire-bonding a p electrode 34 and an n electrode 35 on an upper face of the LED element 33 to a connection terminal 36 on the substrate 31 with a gold wire 37, and then sealing the whole of the LED element 33 with a transparent mold resin 38. However, for the light-emitting device illustrated in FIG. 4, there is the problem that, among the light emitted by the LED element 33, the light having a wavelength of 400 to 500 nm that is output to the upper face side is absorbed by the gold wire, and a part of the light output to the lower face side is absorbed by the die bond adhesive 32, whereby the light-emitting efficiency of the LED element 33 deteriorates.\nConsequently, as illustrated in FIG. 5, flip chip mounting of the LED element 33 has been proposed (Patent Literature 1). In this flip chip mounting technique, a bump 39 is formed on both the p electrode 34 and the n electrode 35. Furthermore, a light-reflecting layer 40 is provided on the surface of the LED element 33 on which the bumps are formed so that the p electrode 34 and the n electrode 35 are insulated. The LED element 33 and the substrate 31 are fixedly connected using an anisotropic conductive paste 41 or an anisotropic conductive film (not illustrated), and curing the paste or film. Consequently, in the light-emitting device illustrated in FIG. 5, the light output toward the upper side of the LED element 33 is not absorbed by the gold wire, and most of the light output toward the lower side is reflected by the light-reflecting layer 40 and is output upwards. As a result, the light-emitting efficiency (light extraction efficiency) does not deteriorate."}
{"text": "The ever expanding use of semiconductor light emitting devices has produced a highly competitive market for these devices. In this market, performance and price are often significant for providing product distinction among vendors.\nSide-emitting LEDs are commonly used in applications requiring a low profile device and a dispersed light pattern, such as backlighted display panels, and thin troffers for general illumination. FIGS. 1A-1C illustrate a conventional technique for providing a side-emitting device, such as disclosed in U.S. Pat. No. 8,109,644, issued 7 Feb. 2012 to Serge Bierhuizen, and incorporated by reference herein.\nFIG. 1A illustrates a substrate 110, upon which semiconductor layers are formed/grown to create multiple light emitting elements 120. The substrate 110 is commonly sapphire, SiC or GaN, and the light emitting element 120 is commonly an active region sandwiched between an n-type semiconductor layer and a p-type semiconductor layer. Pads 130 provide contact to the n and p layers, such that when current flows between these layers, light is emitted from the active region. Gallium nitride (GaN), doped to provide the n-type and p-type semiconductors, is commonly used as the semiconductor forming the light emitting element 120.\nThe substrate 110 may be sliced/diced to provide singular light emitting chips, which may be subsequently placed on another substrate 140, commonly termed a submount, in a ‘flip-chip’ orientation (pads 130 being on the ‘bottom’ of the chip), as illustrated in FIG. 1B. The pads 130 are coupled to contacts 145 on the submount 140 to allow for external electrical connection to the light emitting element 120. With the submount providing structural support to the light emitting element 120, the substrate 110 may be removed.\nTo provide a desired color point (hue and color temperature), a wavelength conversion element 150 is placed upon the light emitting surface 125 of the light emitting element 120, to provide a light emitting device 100 as illustrated in FIG. 1C. The wavelength conversion element may contain one or more wavelength conversion materials, such as phosphors. To provide side emissions, a reflector 160 is situated above the wavelength conversion element 150, so that the light is emitted through the sidewalls 155 of the wavelength conversion element 150. The reflector 160 may be a specular reflector or a diffusing reflector.\nIn some embodiments, light is emitted from all sidewalls of the wavelength conversion element 150; in other embodiments, one or more of the sidewalls of the wavelength conversion element 150 may be reflective, thereby causing the light to be emitted from the remaining non-reflective sidewall(s) 155. Optionally, one or more optical elements, such as collimators, may be used to further direct the side-emitted light in a desired direction.\nGenerally, the wavelength conversion element 150 is attached to the light emitting element 120 soon after the substrate 110 is removed, to protect the light emitting surface 115. Accordingly, the choice of a wavelength conversion element 150 having particular characteristics is generally made prior to the manufacture of the device 100."}
{"text": "The present invention relates generally to poly(imide)s and, more particularly, to high molecular weight poly(imide)s and methods of synthesis thereof.\nSpecialized applications of poly(imide)s capitalize on their unique chemical and physical properties, such as chemical resistance, selective gas permeability, good thermo-oxidative stability, excellent adhesion to metals, high heat performance, and mechanical strength. Poly(imide)s are particularly useful in certain specialized applications, for example in textiles, membranes with selective gas permeability, and in aerospace, aircraft, automotive, and related industries. Manufacture of certain of these articles (particularly selective gas permeable membranes) depends on the solubility and solution properties of the poly(imide). Accordingly, manufacture is preferably being based on solutions having high viscosities and low solid content. Although poly(imide)s having molecular weights of up to 50,000 Daltons are presently used, poly(imide)s having very high molecular weights, (especially greater than 100,000 Daltons) would be even more desirable, as high molecular weight poly(imide)s have high intrinsic viscosities, generally greater than about 0.75 deciliters per gram in the condensed (molten) phase. However, because of these very high viscosities, synthesis and isolation of high molecular weight poly(imide)s using techniques and equipment developed for use with lower molecular weight poly(imide)s is prohibitively difficult.\nThe above-described drawbacks and disadvantages are overcome by a method for the synthesis of a high molecular weight poly(imide) comprising coupling a first poly(imide) precursor having first a functional group and a weight average molecular weight less than about 50,000 Dalton with a second poly(imide) precursor having a second functional group and a weight average molecular weight less than about 50,000 Daltons to form a high molecular weight poly(imide) having a weight average molecular weight greater than 50,000 Daltons.\nIn an alternate embodiment, a method for the synthesis of a high molecular weight poly (imide) comprises coupling an amine-functionalized precursor poly(imide) having a weight average molecular weight of less than about 50,000 Daltons with an anhydride-functionalized precursor poly(imide) having a weight average molecular weight of less than about 50,000 Daltons to provide a poly(imide) having a molecular weight of greater than 50,000 Daltons. The above-described and other features and advantages will be appreciated and understood by those skilled in the art from the following detailed description and appended claims.\nHigh molecular weight thermoplastic poly(imide)s are available from the coupling of lower molecular weight functionalized precursor poly(imide)s. As used herein,xe2x80x9chigh molecular weight poly(imide)sxe2x80x9d refers to poly(imide)s having weight average molecular weights (MW) of greater than 50,000 Daltons, as measured by gel permeation chromatography/laser light scattering/differential viscometry. The method comprises coupling precursor poly(imide)s having complementary functional groups having a weight average molecular weight of less than about 50,000 Daltons. Preferably the method comprises coupling an amine-functionalized precursor poly(imide) having a weight average molecular weight of less than about 50,000 Daltons with an anhydride-functionalized precursor poly(imide) having a weight average molecular weight of less than about 50,000 Daltons. The resultant poly(imide)s have weight average molecular weights greater than 50,000 Daltons, preferably greater than about 60,000 Daltons, more preferably greater than about 75,000 Daltons, even more preferably greater than about 100,000 Daltons, and most preferably greater than about 120,000 Daltons.\nThe functionalized precursor poly(imide)s as well as the high molecular weight poly (imide) are characterized by the presence of an imide group having the general formula xe2x80x94C (O)NRC(O)xe2x80x94. The imide group can be part of either an acyclic or cyclic system within the polymer. Preferred classes of precursor poly(imide)s are aromatic poly(imide)s, characterized by the presence of both an aromatic group and an imide group; poly(amidimide)s, characterized by the presence of both an imide group and an amide group (xe2x80x94C(O)NHxe2x80x94); and poly(etherimide)s, characterized by the presence of both an ether group and an imide group. Useful precursor poly(imide)s include those known in the art which are melt processable, such as those disclosed in U.S. Pat. Nos. 3,803,085 and 3,905,942.\nUseful precursor polyimides have the general formula (I) \nwherein V is a tetravalent linker without limitation, as long as the linker does not impede synthesis or use of the polyimide, and a is more than 1, typically from about 4 to about 1000 or more, and more preferably from about 4 to about 200. Suitable linkers include but are not limited to: (a) substituted or unsubstituted, saturated, unsaturated or aromatic monocyclic and polycyclic groups having about 5 to about 50 carbon atoms, (b) substituted or unsubstituted, linear or branched, saturated or unsaturated alkyl groups having 1 to about 30 carbon atoms; or combinations thereof. Suitable substitutions and/or linkers include, but are not limited to, ethers, epoxides, amides, esters, and combinations thereof. Preferred linkers include but are not limited to tetravalent aromatic radicals of formula (II), such as \nwherein W is a divalent moiety selected from the group consisting of xe2x80x94Oxe2x80x94, xe2x80x94Sxe2x80x94, xe2x80x94C(O)xe2x80x94, xe2x80x94SO2xe2x80x94, Cy H2yxe2x80x94 (y being an integer from 1 to 5), and halogenated derivatives thereof, including perfluoroalkylene groups, or a group of the formula xe2x80x94Oxe2x80x94Zxe2x80x94Oxe2x80x94 wherein the divalent bonds of the xe2x80x94Oxe2x80x94 or the xe2x80x94Oxe2x80x94Zxe2x80x94Oxe2x80x94group are in the 3,3xe2x80x2, 3,4xe2x80x2, 4,3xe2x80x2, or the 4,4xe2x80x2 positions, and wherein Z includes, but is not limited, to divalent radicals of formula (III). \nR in formula (I) includes but is not limited to substituted or unsubstituted divalent organic radicals such as: (a) aromatic hydrocarbon radicals having about 2 to about 20 carbon atoms and halogenated derivatives thereof; (b) straight or branched chain alkylene radicals having about 2 to about 20 carbon atoms; (c) cycloalkylene radicals having about 3 to about 20 carbon atoms, or (d) divalent radicals of the general formula (IV) \nwherein Q includes but is not limited to divalent a divalent moiety selected from the group consisting of xe2x80x94Oxe2x80x94, xe2x80x94Sxe2x80x94, xe2x80x94C(O)xe2x80x94, xe2x80x94SO2xe2x80x94, Cy H2yxe2x80x94 (y being an integer from 1 to 5), and halogenated derivatives thereof, including perfluoroalkylene groups.\nPreferred poly(etherimide) precursors resins comprise more than 1, typically about 4 to about 85 or more, and more preferably about 4 to about 75 structural units, of the formula (V) \nwherein T is xe2x80x94Oxe2x80x94 or a group of the formula xe2x80x94Oxe2x80x94Zxe2x80x94Oxe2x80x94wherein the divalent bonds of the xe2x80x94Oxe2x80x94 or the xe2x80x94Oxe2x80x94Zxe2x80x94Oxe2x80x94 group are in the 3,3xe2x80x2, 3,4xe2x80x2, 4,3xe2x80x2, or the 4,4xe2x80x2 positions, and wherein Z includes, but is not limited, to divalent radicals of formula (III) as defined above.\nIn one embodiment, the polyetherimide may be a copolymer which, in addition to the etherimide units described above, further contains polyimide structural units of the formula (VI) \nwherein R is as previously defined for formula (I) and M includes, but is not limited to, radicals of formula (VII). \nThe polyetherimide can be prepared by any of the methods well known to those skilled in the art, including the reaction of an aromatic bis(ether anhydride) of the formula (VIII) \nwith an organic diamine of the formula (IX) H2 Nxe2x80x94Rxe2x80x94NH2 wherein T and R are defined as described above in formulas (I) and (IV).\nExamples of specific aromatic bis(ether anhydride)s and organic diamines are disclosed, for example, in U.S. Pat. Nos. 3,972,902 and 4,455,410. Illustrative examples of aromatic bis(ether anhydride)s of formula (VII) include: 2,2-bis[4-(3,4-dicarboxyphenoxy)phenyl] propane dianhydride; 4,4xe2x80x2-bis(3,4-dicarboxyphenoxy)diphenyl ether dianhydride; 4,4xe2x80x2-bis (3,4-dicarboxyphenoxy)diphenyl sulfide dianhydride; 4,4xe2x80x2-bis(3,4-dicarboxyphenoxy) benzophenone dianhydride; 4,4xe2x80x2-bis(3,4-dicarboxyphenoxy)diphenyl sulfone dianhydride; 2,2-bis[4-(2,3-dicarboxyphenoxy)phenyl]propane dianhydride; 4,4xe2x80x2-bis(2,3-dicarboxyphenoxy)diphenyl ether dianhydride; 4,4xe2x80x2-bis(2,3-dicarboxyphenoxy)diphenyl sulfide dianhydride; 4,4xe2x80x2-bis(2,3-dicarboxyphenoxy)benzophenone dianhydride; 4,4xe2x80x2-bis (2,3-dicarboxyphenoxy)diphenyl sulfone dianhydride; 4-(2,3-dicarboxyphenoxy)-4xe2x80x2-(3,4-dicarboxyphenoxy)diphenyl-2,2-propane dianhydride; 4-(2,3-dicarboxyphenoxy)-4xe2x80x2-(3,4-dicarboxyphenoxy)diphenyl ether dianhydride; 4-(2,3-dicarboxyphenoxy)-4xe2x80x2-(3,4-dicarboxyphenoxy)diphenyl sulfide dianhydride; 4-(2,3-dicarboxyphenoxy)-4xe2x80x2-(3,4-dicarboxyphenoxy)benzophenone dianhydride and 4-(2,3-dicarboxyphenoxy)-4xe2x80x2-(3,4-dicarboxyphenoxy)diphenyl sulfone dianhydride, as well as various mixtures thereof.\nThe bis(ether anhydride)s can be prepared by the hydrolysis, followed by dehydration, of the reaction product of a nitro substituted phenyl dinitrile with a metal salt of dihydric phenol compound in the presence of a dipolar, aprotic solvent. A preferred class of aromatic bis(ether anhydride)s included by formula (VIII) above includes, but is not limited to, compounds wherein T is of the formula (X) \nand the ether linkages, for example, are preferably in the 3,3xe2x80x2, 3,4 xe2x80x2, 4,3xe2x80x2, or 4,4xe2x80x2 positions, and mixtures thereof, and where Q is as defined above.\nAny diamino compound may be employed in the method of this invention. Examples of suitable compounds are ethylenediamine, propylenediamine, trimethylenediamine, diethylenetriamine, triethylenetetramine, hexamethylenediamine, heptamethylenediamine, octamethylenediamine, nonamethylenediamine, decamethylenediamine, 1,1 2-dodecanediamine, 1,1 8-octadecanediamine, 3-methylheptamethylenediamine, 4,4-dimethylheptamethylenediamine, 4-methylnonamethylenediamine, 5-methylnonamethylenediamine, 2,5-dimethylhexamethylenediamine, 2,5-dimethylheptamethylenediamine, 2,2-dimethylpropylenediamine, N-methyl-bis (3-aminopropyl) amine, 3-methoxyhexamethylenediamine, 1,2-bis(3-aminopropoxy) ethane, bis(3-aminopropyl) sulfide, 1,4-cyclohexanediamine, bis-(4-aminocyclohexyl) methane, m-phenylenediamine, p-phenylenediamine, 2,4-diaminotoluene, 2,6-diaminotoluene, m-xylylenediamine, p-xylylenediamine, 2-methyl-4,6-diethyl-1,3-phenylene-diamine, 5-methyl-4,6-diethyl-1,3-phenylene-diamine, benzidine, 3,3xe2x80x2-dimethylbenzidine, 3,3xe2x80x2dimethoxybenzidine, 1,5-diaminonaphthalene, bis(4-aminophenyl) methane, bis(2-chloro-4-amino-3,5-diethylphenyl) methane, bis(4-aminophenyl) propane, 2,4-bis(b-amino-t-butyl) toluene, bis(p-b-amino-t-butylphenyl) ether, bis(p-b-methyl-o-aminophenyl) benzene, bis(p-b-methyl-o-aminopentyl) benzene, 1,3-diamino-4-isopropylbenzene, bis(4-aminophenyl) sulfide, bis (4-aminophenyl) sulfone, bis(4-aminophenyl) ether and 1,3-bis(3-aminopropyl) tetramethyidisiloxane. Mixtures of these compounds may also be present. The preferred diamino compounds are aromatic diamines, especially m- and p-phenylenediamine and mixtures thereof.\nIn a particularly preferred embodiment, the polyetherimide comprises structural units according to formula (IV) wherein each R is independently p-phenylene or m-phenylene or a mixture thereof and T is a divalent radical of the formula (XI) \nIncluded among the many methods of making the polyimides, particularly polyetherimides, are those disclosed in U.S. Pat. Nos. 3,847,867, 3,814,869, 3,850,885, 3,852,242, 3,855,178, 3,983,093, and 4,443,591. These patents are presented here for the purpose of teaching, by way of illustration, general and specific methods for preparing polyimides.\nIn general, the reactions can be carried out employing well-known solvents, e.g., o-dichlorobenzene, m-cresol/toluene and the like, to effect a reaction between the anhydride of formula (VIII) and the diamine of formula (IX), at temperatures of about 100xc2x0 C. to about 250xc2x0 C. Alternatively, the polyetherimide can be prepared by melt polymerization of aromatic bis(ether anhydride)s (VIII) and diamines (IX) by heating a mixture of the starting materials to elevated temperatures with concurrent stirring. Generally, melt polymerizations employ temperatures of about 200xc2x0 C. to about 400xc2x0 C. Chain stoppers and branching agents may also be employed in the reaction. When polyetherimide/polyimide copolymers are employed, a dianhydride, such as pyromellitic anhydride, is used in combination with the bis(ether anhydride).\nPreferably, the precursor poly(imide)s have a weight average molecular weight (MW ) of from 3,400 to 49,000 grams per mole (xe2x80x9cg/molexe2x80x9d), as measured by gel permeation chromatography/laser light scattering. Such poly(imide)s typically have intrinsic viscosities greater than about 0.2 deciliters per gram, preferably about 0.35 to about 0.7 deciliters per gram measured in m-cresol at 25xc2x0 C.\nIn order to function as precursor poly(imide)s, the above-described polymers must be functionalized with complementary functional groups, that is, functional groups which will allow reaction between the precursor polyimides to form the higher molecular weight products. The above described polymers may be functionalized in one of two ways depending upon the desired functional group. If the desired functional group is the anhydride or amine used to make the polymer itself then producing the precursor poly (imide) merely requires including a molar excess of anhydride or amine in the polymerization reaction. If the desired functional group is other than the anhydride or amine employed in the polymer synthesis then an end-capping agent must employed. Use of endcapping agents is well known in the art. End-capping agents generally comprise either an anhydride or amine functionality in addition to the desired functional group. Choice of the end-capping agent may affect the resulting properties of the high molecular weight poly (imide). Useful complementary functional groups include all functional groups capable of reacting with each other so as to form at least one bond between the precursors thereby coupling the precursor and resulting in a high molecular weight poly(imide). Useful sets of complementary functional groups include, but are not limited to, esters and aliphatic amines, carboxylic acids and amines, acid chlorides and amines, anhydrides and amines, electron rich dienes and dienophiles, and aromatic ring systems containing cyclobutane and functional groups containing a triple bond. A preferred set of complementary functional groups is amines and anhydrides.\nEsters useful as a complementary functional group include all esters having four or less carbons, namely methyl, ethyl, propyl and butyl esters. As is readily understood by one of ordinary skill in the art, methyl esters will react more quickly with a given aliphatic amine than a larger ester such as butyl ester. Additionally, the coupling of an ester with an amine produces the alcohol corresponding to the ester as a side product. The lower alcohols, such as methanol, have a lower boiling point and are easier to remove.\nAmines useful as a complementary functional group with esters are aliphatic and may be primary or secondary aliphatic amines. There are no particular limitations on the identity of the aliphatic amine.\nPrecursors poly(imide)s with carboxylic acid or acid chloride functional groups may be obtained from poly(imide)s with ester functionality by any method known in the art. Precursor poly(imide)s with carboxylic acid or acid chloride functional groups may be coupled with precursor poly(imide)s with amine functional groups. In this case the amine may be a primary or secondary aliphatic amine or an aromatic amine.\nPrecursor poly(imide)s with electron rich diene functional groups may be coupled with precursor poly(imide)s with dienophile type functional groups utilizing the Diels-Alder reaction. Useful electron rich dienes include, but are not limited to, anthracene and furfuryl. Useful dienophiles include but are not limited to maleic anhydride and electron poor alkenes and alkynes. Notably, the coupling of precursor poly(imide)s using a Diels-Alder reaction does not yield any side products.\nPrecursor poly(imide)s comprising a functional group containing an aromatic ring system and cyclobutane may be coupled with a precursor poly(imide) comprising a functional group containing a triple bond which may be terminal or non-terminal. Useful aromatic ring systems with cyclobutane are those in which the cyclobutane is fused between two aromatic rings. Aromatic rings are herein defined as including heteroaromatic rings and ring systems such as pyridine, furan, pyrrole, thiophene, pyrazole, quinoline and fused ring systems such as naphthalene and anthracene. The cyclobutane ring, a strained ring system, reacts with the triple bond to form a six membered aromatic ring without the production of side products.\nIn a preferred embodiment a first precursor poly(imide) comprises at least one amine group, preferably a terminal (end-capping) amine group. Exemplary amine functionalized precursor polyetherimides may be prepared from the reaction of an aromatic bis(ether anhydride) with an organic diamine in which the diamine is present in a molar excess of at least about 2.0 mol%. Preferred amines for synthesis of amine-endcapped poly(etherimide) include, but are not limited to, aliphatic diamines (particularly hexamethylene diamine), gamma-aminopropyl polydimethylsiloxane, m-phenylenediamine, p-phenylenediamine, diethyltoluene diamine or mixtures comprising at least one of the foregoing amines.\nA second precursor poly(imide) comprises anhydride groups, preferably terminal anhydride groups. Poly(imide)s endcapped with anhydride groups are available by synthesis of poly(imide)s as described above, using a molar excess of anhydride relative to amine of at least about 2.0 mol%. Alternative methods for synthesis of anhydride-endcapped poly (imides) are also known, being disclosed, for example, in U.S. Pat. Nos. 3,972,902 and 4,011,198. Preferred anhydrides for the synthesis of anhydride-endcapped poly(imide)s include, but are not limited to, aromatic dianhydrides, fluorinated aromatic dianhydrides, pyromellitic dianhydride, oxydiphthalic anhydride, bisphenol-A dianhydride, and mixtures comprising at least one of the foregoing anhydrides.\nThe degree of endcapping is adjustable by varying amount of excess amine or anhydride. Such amounts are readily determined by one of ordinary skill in the art. Molar excesses in the range from about 2.0 mol% to about 6.0 mol% (based on moles of the other reactant) are generally effective. Under these conditions the poly(imide) resin has between 95 and 150 micro equivalents per gram of acid titratable groups (in the case of amine endcapping, using titration of a chloroform solution with a glacial acetic acid solution of 33% hydrobromic acid ), or base titratable groups (in the case of anhydride endcapping, using titration of a chloroform solution with a 0.1 molar sodium ethoxide solution in ethanol).\nCoupling of the precursors in solution is preferred, as it decreases the viscosity of the system and permits uniform precursor mixing thereby resulting in a more uniform reaction. Use of an appropriate solvent can also facilitate the coupling reaction itself. Suitable solvents are non-reactive, and capable of dissolving the precursor poly(imide)s. Exemplary solvents for the reaction of amine- and anhydride-functionalized poly(etherimide)s include but are not limited to o-dichlorobenzene, trichlorobenzene, N-methylpiperazine (NMP), and mixtures thereof, with o-dichlorobenzene and N-methylpiperazine being preferred. Where water is produced as a reaction by-product, coupling is preferably performed at temperatures above about 100xc2x0 C., preferably in the range of about 100xc2x0 C. to about 350xc2x0 C., most preferably in the range of about 100xc2x0 C. to about 180xc2x0 C. to facilitate the removal of the formed water. However, the reaction temperature preferably is below the boiling point of the solvent. Alternately, chemical drying agents including anhydrides such as acetic anhydride and carbodiimides such as dicyclohexylcarbodiimide can also be used to remove the by-product water. Progress of the reaction may be monitored by measurement of the change in molecular weight and/or intrinsic viscosity (IV).\nCoupling may also be performed using melt-mixing techniques, although solution coupling is preferred, because the high molecular weight products made from the precursors have a high viscosity, which renders melt-mixing difficult.\nThe molecular weight of the product high molecular weight poly(imide) may be determined by direct measurement of the molecular weight. Molecular weight may also be inferred from measurement of the intrinsic viscosity, or, where the percent of endcapping of at least one of the precursors is known, by monitoring the change in concentration of the endcapping by infrared spectroscopy or titration, for example.\nThe high molecular weight poly(imide) may be isolated by conventional techniques, such as solvent evaporation or anti-solvent precipitation. However, isolation of the high molecular weight poly(imide)s, as well as fabrication of articles from the pure, molten poly(imide) can be difficult using present techniques and equipment, which have been developed for use with lower molecular weight polymers. It is therefore preferred to use the formed poly(imide) directly, delaying removal of the solvent by evaporation or anti-solvent precipitation until during and/or after fabrication of the final article.\nAll patents cited are incorporated herein by reference.\nThe following examples are given by way of illustration, and not by way of limitation."}
{"text": "In medicine, mechanical ventilation is a method to mechanically assist or replace spontaneous breathing of a patient using a machine called a ventilator. Mechanical ventilation is often a life-saving intervention, but carries many potential complications including pneumothorax, airway injury, alveolar damage, and/or ventilator-associated pneumonia, thereby requiring that a respiratory care practitioner operate the ventilator.\nThe delivery of a “gold standard” of care in mechanical ventilation to a patient population relies on having a sufficient number of acute care ventilators on hand as well as a requisite number of respiratory care practitioners to properly operate them. Recognizing that in a severe pandemic like the 1918 Spanish Flu pandemic, or a mass casualty event, such as a major earthquake, hurricane, or terrorist incident, the very real possibility exists that patient loads generated by such events will initially exceed the number of stockpiled ventilators and/or the requisite number of respiratory care practitioners on hand to provide even a modified gold standard treatment to a large number of patients requiring mechanical ventilation. Even if a sufficient number of mechanical ventilators are stockpiled in a particular area, the large number of immediate casualties in such a catastrophic event could overwhelm the limited number of respiratory care practitioners available to provide the necessary expertise to operate all of the ventilators required to treat a large number of patients, especially those patients not being treated by emergency personnel or at a healthcare facility, such as a hospital. Accordingly, there is a need for a ventilation system that may be operated by a non-expert with no or little experience in the operation of a ventilator as well as a ventilator adapted for seamless adjustment of ventilator functions and parameters by a respiratory care practitioner. There is also a need for a ventilation system having various modes of operation that can be utilized to mechanically ventilate patients outside a hospital setting without an AC power source."}
{"text": "The compressor art, namely, such machinery or devices as used to compress gases, is old and differing structural designs are known, many of which are patented.\nThis invention is concerned with a type of air compressor incorporating a globoid-worm having a plurality of threads in meshing relation with at least one toothed pinion, as exemplified by U.S. Pat. No. 3,945,778, Mar. 23, 1976, and U.S. Pat. No. 3,804,564, Apr. 16, 1954, both granted to Bernard Zimmern. More specifically, the subject invention represents an improvement in such Zimmern compressors insofar as manufacturing costs are reduced, efficiency is improved, and reliability is increased, all as will become evident from the detailed description of the invention set forth hereinafter."}
{"text": "Radar signals have been commonly used by police for some time to determine the speed of motor vehicles. In response to radar speed monitoring and to signal motor vehicle operators when such monitoring is taking place, police radar detectors have likewise been used for almost a coincident period of time. Currently available radar detectors indicate the presence of radar signals, the frequency band of detected signals, the direction from which the radar signals originate and the relative field strength of detected signals. In addition, the current radar detectors can also display information about their current mode of operation and the number of detected radar signals at any given time. The widely varying operating procedures for using police radar and the proliferation of other signals assigned to the same frequency bands as police radar has led to the need for police radar detectors which give more information than has been contemplated in the past.\nFor example, police radar units are often deployed along the side of the roadway, to measure a motor vehicle's speed as it advances toward the unit. In some instances, the police radar may measure the speed of a motor vehicle after it has passed and is travelling away from the unit. Police vehicles may also be equipped with radar units which are operated while the police vehicles are moving, using reflections from stationary objects to measure the speed of the police vehicle itself, and reflections from a target vehicle to measure relative speed, and using both measured speeds to determine the actual speed of the target vehicle.\nWhen a vehicle equipped with a radar detector is within microwave range of a police radar unit which transmits a radar signal, the detector alerts the vehicle operator that the signal is present typically through a combination of audible and visual signals. As the user approaches the source of the radar signal, signal strength of the detected radar signal typically increases. As the vehicle passes the source of the radar signal, the indicated signal strength drops, usually very quickly, since a forwardly aimed directional antenna of the detector is no longer pointed in the general direction of the signal source, and is now responding to reflections of the signal from objects in front of the vehicle. If the police radar unit is at the side of the road in a clearly visible location, the vehicle operator can easily correlate the detected signal and its apparent source.\nOften, a driver using a radar detector may be travelling at a relatively high rate of speed and would benefit from receiving all the available information about a detected radar signal as quickly and concisely as possible. Thus, there remains a need for police radar detectors that efficiently and effectively convey as much information to a driver as easily and quickly as possible."}
{"text": "This invention relates generally to local area networks (LANs) of the Ethernet type and, more particularly, to techniques for improving fairness of channel access in an Ethernet LAN. Ethernet is a commonly used name for a LAN that uses a network access protocol referred to as Carrier Sense Multiple Access with Collision Detection (CSMA/CD). The CSMA/CD protocol is defined in ANSI/IEEE Std 802.3, published by the Institute of Electrical and Electronics Engineers, Inc., 345 East 45th Street, New York, N.Y. 10017, and referred to in this specification as \"the IEEE 802.3 standard.\" The standard is for a 10 Mbps (megabits/sec) CSMA/CD channel, but it will be understood that the present invention may also be applicable to a 100 Mbps channel.\nUnder the CSMA/CD rules for access to a network bus or cable, which will be referred to as the channel, any node or station wishing to transmit must first \"listen\" to make sure that the channel is clear before beginning to transmit. All nodes on the network have equal priority of access and may begin transmitting as soon as the channel is clear and a required inter-packet delay of 9.6 .mu.s (microseconds) has elapsed. However, if a first node that has started transmitting detects a \"collision\" with a transmission from another node, the first node continues transmitting for a short time to make sure that all nodes wishing to transmit will detect the collision. Every other node detecting the collision also continues to transmit for a short time. Then each node that has detected a collision terminates transmission of the packet or frame. The nodes involved in the collision wait for the required interpacket delay of 9.6 .mu.s and then select random, and therefore usually different, delay times, referred to as backoff times, before trying transmission of the same packet again.\nMore specifically, the IEEE 802.3 standard defines a collision backoff procedure referred to as \"truncated binary exponential backoff.\" When a transmission attempt has terminated due to a collision, it is retried by the transmitting node, after a selected backoff time, until either the transmission is successful or a maximum number of attempts have been made and all have terminated due to collisions. The backoff time is selected by each node as an integral multiple of the \"slot time,\" which is the maximum round-trip propagation time for the network, i.e. the time to propagate a data packet from one end of the network to the other, and back. The slot time defined by the IEEE 802.3 standard is 51.2 .mu.s. The number of slot times selected as the backoff time before the nth retransmission is chosen as a randomly distributed integer r in the range: EQU 0.ltoreq.r&lt;2.sup.k, where k=min (n, 10).\nTherefore, for the first attempted retransmission the backoff time is selected as 0 or 1 slot times, for the second attempted retransmission the backoff time is selected as 0, 1, 2 or 3 slot times, for the third attempted retransmission the backoff time is selected as 0, 1, 2, 3, 4, 5, 6 or 7 slot times, and so forth. The maximum backoff time, for the tenth attempted retransmission may be up to 2.sup.10 -1, or 1,023 slot times, i.e. 52.4 ms.\nThe IEEE 802.3 standard is designed to achieve optimal performance, in throughput, latency and fairness, when the number of nodes in the network is relatively large. When the number of active nodes, i.e. nodes having data packets to transmit, is small, e.g. two nodes, the IEEE 802.3 standard exhibits an undesirable effect referred to as the capture effect, in which one of the nodes may effectively capture the channel and transmit a succession of data packets in a back-to-back fashion for a relatively long period of time. The capture effect is best understood from a simple example.\nSuppose there are two active nodes, A and B, in a network and they begin transmitting at approximately the same time, resulting in a first collision. They each select backoff times of 0 or 1 slot time, in accordance with the standard backoff algorithm. Suppose further that node A selects zero backoff time and node B selects one slot time as its backoff time. Obviously, node A will be able to successfully transmit on its first attempt at retransmission. Node B will wait a full slot time before making its first retransmission attempt, but by this time node A has started transmitting a second data packet, i.e. node A is transmitting successive packets back-to-back, with only the required interpacket gap (IPG) separating them. Node B experiences a second collision on its first data packet, but for node A it is still a \"first\" collision for trying to transmit a second data packet. In accordance with the standard backoff algorithm, node A selects a backoff of 0 or 1 slot time, but node B selects a backoff of 0, 1, 2 or 3 slot times. Therefore, node B has only one chance in eight of winning access to the channel (if node B chooses 0 and node A chooses 1). Node A, however, has five chances in eight of winning channel access. (Two of the eight possibilities will result in another collision.) Clearly, it is more probable, actually five times more probable, that node A will retain access to the channel, as compared with node B. If the probable occurs, and node A transmits its second data packet,the situation becomes progressively more difficult for node B to transmit. If another collision occurs on node B's next attempt to retransmit, this will be node B's third attempt, but node A's \"first\" again. Node A's backoff is still 0 or 1 slot time, but node B's is selected from 0 through 7 slot times. Node A is then thirteen times as likely to win access to the channel on the next attempted retransmission. For the nth attempt at retransmission by node B, node A will be (2.sup.n -3) times more likely to win access to the channel. For the maximum value of n, i.e. 10, the unfairness factor between the two active nodes will be 2.sup.10 -3, or 1,021.\nFor network protocols that require a receiving node to acknowledge that it has received a packet, the situation is further aggravated by the need to transmit acknowledgment messages on the network. For example node A captures the channel and is sending successive packets to node B, node B may be unable to acknowledge even the first received packet. Node A will eventually have to give up the channel to allow acknowledgements to flow back from node B, but this is an inefficient process and channel utilization goes down.\nIt will be appreciated from the foregoing that there is a need for improvement over the standard approach for selecting random backoff times in a network using the CSMA/CD protocol. Ideally, any improved technique should still comport with the objects of the IEEE 802.3 standard and should be compatible with nodes that use the standard without improvement. The present invention achieves these ends, as will become apparent from the following summary.\nOne approach has been to assign special backoff times to a node based on the node's recent transmission attempt history. These special backoff algorithms, however, sometimes have the undesirable effect of reducing channel throughput."}
{"text": "1. Technical Field\nThe illustrative embodiments relate to a method and apparatus for automatically provisioning an emergency phone call placed by a vehicle computing system (including, but not limited to, recognizing local emergency call protocol, switching languages to a local language, etc.).\n2. Background Art\nMany vehicles come equipped with a variety of automatic computing systems. As computers grow more powerful, it is possible to add a multitude of functions to these computing systems.\nFor example, without limitation, the FORD SYNC system gives a user access to a powerful set of tools when traveling. The user can receive and make phone calls through the system, check email, send and receive text messages, check sports scores and weather, order carry out/drive-through food, play games, receive driving directions, etc.\nIn these illustrative examples, one method of connection is done through a user's cellular phone or other nomadic device. The vehicle computing system, such as, but not limited to, the FORD SYNC system, connects to the user's nomadic device, and uses that device to establish a connection with a remote network. Using connection options like Voice Over IP (VOIP) and/or a data connection, information can be sent to and from the vehicle computing system.\nIn addition to user initiated communication, it may be possible for the vehicle computing system to initiate communication on its own. The system may need to check a remote network for updates, or may need to process ongoing data transfer for a previously requested service.\nIn at least one instance, it may be desirable to have the vehicle equipped with functionality for automatically placing emergency phone calls in the event of a vehicle accident. Such a system is described in some detail in pending U.S. application Ser. No. 11/769,346, entitled “METHOD AND SYSTEM FOR EMERGENCY NOTIFICATION”, filed Jun. 27, 2007; and U.S. application Ser. No. 12/399,513, entitled “METHOD AND SYSTEM FOR EMERGENCY CALL HANDLING”, filed Mar. 6, 2009; and U.S. application Ser. No. 12/607,244, filed Oct. 28, 2009, entitled “METHOD AND SYSTEM FOR EMERGENCY CALL PLACEMENT”. The contents of which are incorporated herein by reference.\nTypically, when a call is placed, the language of the call will be determined by the language spoken by the vehicle computing system when communicating with the user, or by the default language as determined for a particular region, or determined at the end line of manufacturing, etc. For example, if the user has the vehicle computing system set to English, then the call will be placed in English. This has the potential, however, to create complications if the user drives the vehicle across a border between countries.\nFor example, if the user drove from America to Mexico, then the language may be set to English, but the emergency call may be placed to a Mexican Emergency Operator, and the operator may only speak Spanish. This could impair or even prevent successful completion of the emergency call."}
{"text": "Couplings and transition fittings of the type having a gripping member are well known, see for example Applicant's U.S. Pat. No. 5,593,186. In general these couplings include a nut, body, gripping member having a plurality of inwardly directed barbed projections, as well as a sub-assembly of a compression sleeve and resilient gasket ring that is typically permanently attached to the inner end of the compression sleeve. The nut and the gripping member include abutting tapered surfaces. A pipe end is freely insertable through the sub-assembly and into a socket in the body of the coupling. As the nut is tightened it exerts both a radial and longitudinal force on the grip ring that in turn longitudinally exerts a force on the sleeve and gasket ring. The gasket ring encounters an abutment in a socket in the body, and further tightening of the nut then further radially compresses the gripping member so that its inwardly directed barbed projections engage the outer surface of a pipe to which the coupling is attached. There may also be provided abutment surfaces within the body of the coupling and the nut that abut inner and outer surfaces of a flange outstanding from the outer end of the compression sleeve.\nOne skilled in the art will appreciate that as the nut is tightened, the tapered surface of the nut acts upon the tapered surface of the gripping member causing it to be both compressed and urged into the body. As the gripping member engages the pipe it is also drawn into the assembly.\nWhilst this works well, couplings are generally designed for a pipe of a particular diameter requiring either the retailer or the end user to stock various couplings. A further problem is that couplings may be manufactured either to meet the metric or the imperial standard.\nThe main object of this invention is therefore to overcome the abovementioned problems or at least provide the public with a useful alternative by providing improvements whereby the coupling can accommodate pipes of two different diameters. In particular the coupling can accommodate two pipes whose diameter may be different by up to 5-10 mm."}
{"text": "Metallic glasses are metallic alloys that have a glassy phase with an amorphous atomic structure in a solid state. The glassy phase is believed to be a metastable phase and not a thermodynamically stable phase in a solid state. As a result, metallic glasses are typically formed by quenching from a liquid state to reduce or even avoid nucleation and growth of crystalline phases during solidification. As a result, casting large articles of metallic glasses may be difficult because large articles may not be quenched at sufficiently high rates."}
{"text": "This invention relates to occupant restraint systems and more particularly to occupant restraint systems of the type which releasably secure the cushion against initial deployment normal to or toward the occupant and direct such initial deployment laterally of the occupant\nU.S. Pat. No. 5,009,452, Occupant Restraint System, issued Apr. 23, 1991 to Gregory A. Miller and assigned to the assignee of this invention, discloses an occupant restraint system which includes a reaction member mounted to the support of the restraint module and located above the inflator of an occupant restraint system. The reaction member is sealed to the upper wall portion of the occupant restraint cushion to prevent the impingement of pressure fluid from the inflator against such upper wall portion during initial deployment of the cushion and thereby restrict such initial deployment laterally of the occupant through the unfolding of the folded side and end walls of the cushion.\nThis invention provides an occupant restraint system having a reaction member of the type disclosed in the Miller patent and having additional features.\nOne of the features of the occupant restraint system of this invention is that the reaction member is mounted to the support by a flexible strap or mounting member having a bight which loops through a pair of spaced slots or openings located in the center area of the reaction member. The legs of the strap pass around the inflator and downwardly through the inflator receiving opening of the support to the underside thereof. The ends of the legs are doubled over and secured to the underside of the support. If, desired the legs need not be secured to the support since the doubled over ends prevent the ends of the legs from being pulled upwardly of the support between the inflator and the inflator receiving opening of the support. By mounting the reaction member to the support by a flexible strap, the size of the inflator module can be reduced since the folded end walls of the cushion can be moved inward toward the inflator and occupy the space otherwise occupied by the supports for the reaction member in the Miller patent. Additionally, the mounting member permits the reaction member to shift or move slightly relative to the support and inflator if required.\nAnother feature is that the reaction member is pocketed to the inside of the upper wall portion of the cushion to seal the interface between the reaction member and such upper wall portion of the cushion. In one embodiment of the invention, the pocketing includes flaps of flexible material, such as the same material as the cushion, which are sewn or otherwise secured to the inside of the upper wall portion of the cushion adjacent the periphery of the reaction member and overlap or overlie such periphery. The flaps are of truncated triangular shape to leave the center area of the reaction member uncovered and not interfere with the passage of the strap therethrough. The side edge portions of the flaps overlie or overlap each other. The flaps are releasably secured to the reaction member by a retainer member which sandwiches the flaps to the reaction member. The retainer member is centrally secured to the reaction member and releasably secures the flaps to the reaction member by shearable fasteners which extend through the overlapping side edge portions of the flaps. The retainer member is provided with slots or openings in general alignment with those of the reaction member so that the bight of the mounting member can extend therethrough. In another embodiment of the invention, the pocketing includes a sheet of flexible material of the same general shape as the reaction member but slightly larger. The edge portions of the sheet are sewn or otherwise secured to the inside of the upper wall portion of the cushion adjacent the periphery of the reaction member. The center area of the sheet of flexible material is open and slits, perforations, or other types of tear lines extend diagonally from such center opening toward the corner portions of the sheet of flexible material. The slits, perforations, or other types of tear lines terminate slightly short of the corner portions of the reaction member so that the sheet of flexible material pockets the reaction member to the upper wall portion of the cushion. This embodiment also includes a retainer member centrally secured to the reaction member and sandwiching the sheet of flexible material to the reaction member. However, there is no need for the retainer member to include any shearable fastening means for releasably securing the sheet of flexible material to the reaction member.\nWhen the folded undeployed cushion initially receives pressure fluid from the inflator, the initial deployment of the cushion is laterally of the occupant as the folded side and end walls of the cushion unfold. The pocketing of the reaction member to the upper wall portion of the cushion by the flaps or sheet of flexible material seals the interface between the reaction member and such upper wall portion of the cushion and prevents the entry of the pressure fluid therebetween. Thus, the upper wall portion of the cushion remains generally stationary and does not move toward the occupant during initial deployment of the cushion. As the pressure fluid continues to fill the cushion, the side and end walls unfold into continuations of the pocketed and sealed upper wall portion. When the pressure and volume of the pressure fluid reach certain levels, the pocket is ruptured by (1) the flaps being pulled from between the retainer member and reaction member as the shearable fasteners shear, or (2) the sheet of flexible material tearing or rupturing along the tear lines into flaps connected at the corner portions of such sheet, to permit the flaps to move with the upper wall of the cushion toward the occupant. The flexible mounting member retains the reaction member and retainer member in place. Since the flaps are formed of flexible material, the occupant cannot feel or engage the flaps through the cushion should the occupant engage the outer side of the upper wall of the cushion."}
{"text": "As computing technology has advanced, computers have become increasingly interconnected. This interconnection allows large amounts of data to be available to computers, such as Web pages from numerous Web sites. Having such a large amount of data available is useful, but is not without its problems. One such problem is that it can be time consuming to have the Web pages the user desires obtained and displayed. This can lead to a frustrating user experience due to delays in the user being able to see the Web pages he or she desires."}
{"text": "The invention relates to the field of integrated circuit chip carriers. Chip carriers provide an alternative method of packaging integrated circuit chips to the well-known dual in-line package (DIP).\nKnown forms of chip carrier consist of a usually rectangular base of ceramic material. A central area on the base is provided, on to which an integrated circuit chip or die is mounted and secured, for example using an epoxy resin adhesive. Leads are then connected between the appropriate points on the die to metal bonding pads provided on the carrier base. A preformed metallization pattern on the carrier base connects the bonding pads to terminal contacts formed around the periphery of the base. A metal or ceramic lid is then secured over the base so as to enclose and protect the die. The ceramic chip carrier is therafter mounted on a printed circuit board together with a series of other chip carriers and electronic components that go to form the complete circuit.\nThe use of ceramic material for the chip carrier is expensive and problems can also arise when the ceramic base is later bonded to a conventional epoxy/glass printed circuit board, due to the differential expansion between the two materials. In extreme environmental conditions the ceramic chip carrier can break away from the printed circuit board."}
{"text": "The present invention relates to a honeycomb structure with a sensor insertion hole."}
{"text": "It is well known in the art to grow mushrooms in host wood. Conventionally, this is accomplished by employing wood logs, wood pieces, and even wood in particulate form as the host or growing medium, the wood having been previously removed from a tree and maintained in a controlled environment conducive to mushroom growing.\nSuch prior art approaches are quite expensive and often do not lend themselves to large scale mushroom production. The following U.S. patents have been located and are believed to be representative of the current state of the prior art: U.S. Pat. No. 3,563,461, issued Feb. 16, 1971, U.S. Pat. No. 4,833,821, issued May 30, 1989, U.S. Pat. No. 4,918,859, issued Apr. 24, 1990, U.S. Pat. No. 1,833,089, issued Nov. 24, 1931, U.S. Pat. No. 4,763,440, issued Aug. 16, 1988, U.S. Pat. No. 4,646,465, issued Mar. 3, 1987, U.S. Pat. No. 5,001,858, issued Mar. 26, 1991, and U.S. Pat. No. 5,222,665, issued Jun. 29, 1993."}
{"text": "As a next-generation system of the W-CDMA (Wideband Code Division Multiple Access) system, HSDPA (High Speed Downlink Packet Access) system, HSUPA (High Speed Uplink Packet Access) system and the like, an LTE system has been studied by 3GPP (3rd Generation Partnership Project) which is a standards body of the W-CDMA. In the LTE system as a radio access system, an OFDM (Orthogonal Frequency Division Multiplexing) scheme and an SC-FDMA (Single-Carrier Frequency Division Multiple Access) scheme have been studied as prospective candidate schemes to be applied to the downlink communications system and the uplink communications system, respectively (see, for example, Non Patent Document 1).\nThe OFDM scheme is a multi-carrier transmission scheme in which a frequency band is divided into plural sub-carriers having narrower frequency bands, and data are mapped onto the sub-carriers. By closely and orthogonally arranging the sub-carriers along the frequency axis, faster transmission is expected to be achieved and the efficiency of using the frequency band is also expected to be improved.\nThe SC-FDMA scheme is a single carrier transmission scheme in which a frequency band is divided with respect to user equipment (hereinafter may be referred to as user equipment (UE) terminal or mobile station) in a manner so that different frequencies can be separately used among plural terminals (user equipment terminals); and as a result, interference between terminals may be easily and effectively reduced. Further, preferably, in the SC-FDMA scheme, a range of transmission power fluctuation may be made smaller; therefore lower energy consumption of terminals can be achieved and a wider coverage area can be obtained.\nIn both uplink and downlink of the LTE system, communications can be performed by allocating one or more resource blocks to the mobile station. More specifically, the resource blocks may be shared among plural mobile stations in the system. In the LTE system, the base station (hereinafter may be referred to as base station apparatus) determines which resource blocks are to be allocated to which mobile station from among plural mobile stations with respect to each Sub-frame having 1 ms period (this allocation process may be called Scheduling). The Sub-frame may also be called a TTI (Transmission Timer Interval). In downlink, the base station apparatus transmits a shared channel using one or more resource blocks to the mobile station selected by the Scheduling. The shared channel is called a Physical Downlink Shared Channel (PDSCH). In uplink, the mobile station selected by the Scheduling transmits the shared channel using one or more resource blocks. This shared channel is called a Physical Uplink Shared Channel (PUSCH).\nIn a communication system using the shared channels, it is required to perform (report) Signaling designating which shared channel is allocated to which user equipment terminal. In the LTE system, a control channel is used for the Signaling, and the control channel is called a Physical Downlink Control Channel (PDCCH) or a Downlink L1/L2 Control Channel (DL-L1/L2 Control Channel). The Physical Downlink Control Channel (PDCCH) may include information items such as Downlink Scheduling Information, Acknowledgement information (ACK/NACK), Uplink Scheduling Grant, an Overload Indicator, Transmission Power Control Command Bit and the like (see, for example, Non Patent Document 2).\nThe Downlink Scheduling Information and Uplink Scheduling Grant correspond to the information on which the Signaling is to be performed. The Downlink Scheduling Information includes the information of the shared channel in downlink. More specifically, the Downlink Scheduling Information includes allocation information of the Resource Blocks in downlink, identification information of the user equipment (UE-ID), the number of streams, information of Precoding Vector, data size, modulation scheme, information of Hybrid Automatic Repeat reQuest (HARQ) and the like.\nThe Uplink Scheduling Grant includes information of the shared channel in uplink. More specifically, the Uplink Scheduling Grant includes allocation information of the Resources in uplink, the UE-ID, data size, modulation scheme, information of transmission power in uplink, information of a Demodulation Reference Signal in Uplink MIMO and the like.\nIn downlink, a Common Control Physical Channel (CCPCH) is also transmitted. The Common Control Physical Channel (CCPCH) includes a Broadcast Channel (BCH). The Broadcast Channel (BCH) transmitted in the Common Control Physical Channel (CCPCH) may be called a static Broadcast Channel. In addition to the static Broadcast Channel, there is also provided a dynamic Broadcast Channel (Dynamic part). The dynamic Broadcast Channel is mapped to the Physical Downlink Shared Channel (PDSCH). In this case, the Physical Downlink Shared Channel (PDSCH) transmits the Downlink Scheduling Information for the dynamic Broadcast Channel.\nAn uplink control channel transmits a downlink Channel Quality Indicator (CQI), the Acknowledgement information with respect to the Physical Downlink Shared Channel (PDSCH) and the like. The Channel Quality Indicator (CQI) is used in, for example, the Scheduling process and an Adaptive Modulation and Coding Scheme (AMCS) with respect to the Physical Downlink Shared Channel (PDSCH).\nAs described above, in the proposed mobile communication system, basically, the information indicating which shared channel is allocated to which user equipment terminal with respect to each Sub-frame is required to be reported by Signaling using the Downlink L1/L2 Control Channel (DL-L1/L2 Control Channel). Further, this Scheduling is to be required (reported) regardless of the length of the packet data to be transmitted. Therefore, even when the packet data having a short length are required to be frequently transmitted, it is required to perform Signaling designating which shared channel is allocated to which user equipment terminal with respect to each packet data transmission using the Downlink L1/L2 Control Channel (DL-L1/L2 Control Channel). As a result, the ratio of the radio resources allocated to the control channel increases, the overhead rate may be accordingly increased, and the number of radio resources to be allocated to the data channel may be reduced. Typical examples of the packet data that have a short length and that are generated frequently may be voice packet data, VoIP, real-time data, and the like.\nTo overcome the problems, a method called Persistent Scheduling has been proposed. In this method, a downlink data channel (typically, the voice packet data) is transmitted using a certain fixed transport format based on a specific cycle such as 20 ms. The transport format includes information items necessary to decode the data channel, the information items indicating the modulation scheme, the channel coding rate and the like. For example, the modulation scheme and the channel coding rate are fixed to QPSK and ⅓, respectively, and the data are known to (i.e., shared between) the base station and the user equipment terminals. Therefore, the user equipment terminals may adequately receive the downlink shared channel (DL-SCH) even when the control channel such as the L1/L2 control channel with respect to each Sub-frame is not reported. Such Persistent Scheduling is described in, for example, Non Patent Documents 3 and 4.                Non Patent Document 1: 3GPP TR 25.814 (V7.0.0), “Physical layer Aspects for Evolved UTRA,” June 2006        Non Patent Document 2: 3GPP R1-070103, Downlink L1/L2 Control Signaling Channel Structure: Coding        Non Patent Document 3: R1-051511, 7-11th Nov., 2005, 3GPP TSG-RAN WG1 #43, Qualcomm Europe, pp 2, Section 3        Non Patent Document 4: R2-060550, 13-17 Feb., 2006, 3GPP TSG-RAN WG2 #51, Qualcomm Europe, pp 1, Section 2        "}
{"text": "Dynamic random access memories (DRAM's) are used as main memory in many of today's computer systems. One of the factors that have led to the popularity of the DRAM has been the DRAM's simple cell structure. Because each DRAM storage cell consists of just a single capacitor, it is possible to pack a very large number of storage cells into a very small amount of chip space. Consequently, with DRAM technology, it is possible to manufacture very high-density, low cost memories.\nWith reference to FIG. 1, there is shown a functional diagram of a typical DRAM 100. As shown in FIG. 1, a DRAM 100 comprises a plurality of memory cells 102 arranged in a plurality of rows and columns. Each row of memory cells 102 is coupled to one of the wordlines 104 of the DRAM 100, and each column of memory cells 102 is coupled to one of the bitlines 106. By specifying a wordline 104 and a bitline 106, a particular memory cell 102 can be accessed.\nTo enable access to the various memory cells 102, there is provided a row decoder 112 and a column decoder 114. The row decoder 112 receives a row address on a set of address lines 116, and a row address strobe (RAS) signal on a control line 118, and in response to these signals, the row decoder 112 decodes the row address to select one of the wordlines 104 of the DRAM 100. The selection of one of the wordlines 104 causes the data stored in all of the memory cells 102 coupled to that wordline 104 to be loaded into the sense amplifiers (sense amps) 108. That data, or a portion thereof, may thereafter be placed onto the data bus 110 (for a read operation).\nWhat portion of the data in the sense amps 108 is actually placed onto the data bus 110 in a read operation is determined by the column decoder 114. More specifically, the column decoder 114 receives a column address on the address lines 116, and a column address strobe (CAS) signal on a control line 120, and in response to these signals, the column decoder 114 decodes the column address to select one or more of the bitlines 106 of the DRAM 100. The number of bitlines 106 selected in response to a single column address may differ from implementation to implementation, and is referred to as the base granularity of the DRAM 100. For example, if each column address causes sixty-four bitlines 106 to be selected, then the base granularity of the DRAM 100 is eight bytes. Defined in this manner, the base granularity of the DRAM 100 refers to the amount of data that is read out of or written into the DRAM in response to each column address/CAS signal combination (i.e. each CAS or column command).\nOnce the appropriate bitlines 106 are selected, the data in the sense amps 108 associated with the selected bitlines 106 are loaded onto the data bus 110. Data is thus read out of the DRAM 100. Data may be written into the DRAM 110 in a similar fashion. A point to note here is that in a typical DRAM, the address lines 116 are multiplexed. Thus, the same lines 116 are used to carry both the row and column addresses to the row and column decoders 112, 114, respectively. That being the case, a typical memory access requires at least two steps: (1) sending a row address on the address lines 116, and a RAS on the control line 118; and (2) sending a column address on the address lines 116, and a CAS on the control line 120.\nA timing diagram illustrating the various steps carried out during typical DRAM read cycles is shown in FIG. 2. As shown, to initiate a read cycle, a row address 208(1) is placed onto the address lines 116, and a RAS signal 202(1) is asserted on the RAS control line 118. Then, a column address 210(1) is sent onto the address lines 116, and a CAS signal 204(1) is asserted on the CAS control line 118. A short time thereafter, the data 206(1) stored at the locations indicated by the row address 208(1) and the column address 210(1) appear on the data bus 110. Data is thus extracted from the DRAM 100. After the first set of data 206(1) disappears from the data bus 110, a second read operation may be initiated. Like the first read operation, the second read operation begins with a row address 208(2) on the address lines 116, and a RAS signal 202(2) on the RAS control line 118. Then, a column address 210(2) is sent onto the address lines 116, and a CAS signal 204(2) is asserted on the CAS control line 118. A short time thereafter, the data 206(2) stored at the locations indicated by the row address 208(2) and the column address 210(2) appear on the data bus 110. The second read operation is thus completed. Additional successive reads may be carried out in a similar fashion.\nNotice from the timing diagram of FIG. 2 that, for individual read cycles, there is substantial idle time between successive data sets 206 on the data bus 110. During this idle time, the data bus 110 is not utilized and no data is being transferred. The more idle time there is, the lower the utilization rate of the data bus 110, and the lower the utilization rate, the longer it will take for an external component (such as a CPU) to extract data from the DRAM 100. Since almost all operations of a computer require the use of memory, the longer it takes to get data from a memory, the slower the performance of the overall computer system. Thus, low bus utilization can have a direct negative impact on the overall performance of a computer system.\nTo improve data bus utilization, several techniques have been developed. One such technique involves the use of a “burst” mode of operation. Basically, in burst mode, rather than implementing just one column access for each RAS command, a plurality of column accesses are carried out for each RAS command. This results in consecutively accessing multiple sets of data from the same row of a DRAM 100. A timing diagram illustrating the operation of burst mode is shown in FIG. 3. More specifically, FIG. 3 depicts two burst mode read cycles, with each read cycle being directed to a different row of the DRAM.\nTo initiate a burst mode read cycle, a row address 308(1) is placed onto the address lines 116, and a RAS signal 302(1) is asserted on the RAS control line 118. Then, a column address 310(1) is sent onto the address lines 116, and a CAS signal 304(1) is asserted on the CAS control line 118. In response to the column address 310(1) and the CAS signal 304(1), the DRAM internally generates a plurality of additional column addresses. These additional column addresses are generated based upon the column address 310(1) that is provided, and a predetermined scheme. For example, the additional column addresses may be generated by incrementing the provided column address 310(1), decrementing the column address 310(1), or by manipulating the column address 310(1) in some other manner. The number of additional column addresses generated by the DRAM depends upon the burst length that the DRAM is implementing. In the example shown in FIG. 3, the burst length is four; thus, three additional column addresses are generated by the DRAM.\nAs the provided column address 310(1) is received, and as each additional column address is generated, they are applied by the DRAM to access a particular set of data. These addresses are applied in succession so that multiple sets of data are accessed from the same row of the DRAM. A short time after the application of each column address, data 306 stored at the locations indicated by the row address 308(1) and the applied column address starts to appear on the data bus 110. Because this data 306 is extracted from the DRAM 100 in response to consecutive applications of column addresses, there is no idle time between the sets of data 306(1)-306(4) on the data bus 110. As a result, data bus utilization is improved.\nAfter the first set of data 306(1)-306(4) disappears from the data bus 110, a second burst mode read operation may be initiated. Like the first read operation, the second read operation begins with a row address 308(2) on the address lines 116, and a RAS signal 302(2) on the RAS control line 118. Then, a column address 310(2) is sent onto the address lines 116, and a CAS signal 304(2) is asserted on the CAS control line 118. In response, the DRAM generates three additional column addresses, and applies the provided column address 310(2) and the additional column addresses in succession to access multiples set of data from the same row. Shortly after each column address is applied, data 306 stored at the locations indicated by the row address 308(2) and the applied column address appears on the data bus 110. Because this data 306 is extracted from the DRAM 100 in response to consecutive applications of column addresses, there is again no idle time between the sets of data 306(5)-306(8) on the data bus 110. The second read operation is thus completed. Additional successive reads may be carried out in a similar fashion.\nSeveral aspects of burst mode operation should be noted. First, notice that burst mode significantly increases output data granularity. More specifically, because a burst mode memory request involves multiple column accesses, the data extracted from the DRAM in response to a burst mode request is not just one base granularity in size, but rather is a multiple of the base granularity, where the multiple is equal to the burst length. Thus, in the example shown in FIG. 3, the output data granularity of the DRAM is four times the base granularity. This may pose a problem in some implementations. For example, in some applications, a CPU may wish to access only one base granularity of data at a time. If burst mode is implemented in such an application, then all of the data after the first granularity will be dropped by the CPU. In such a case, even though data bus utilization is improved by the use of burst mode, overall system efficiency is not improved because the extra data from the memory is not used. From an efficiency point of view, the end result is the same as if burst mode were not implemented at all. In such applications, burst mode does not provide a useful solution.\nA second aspect to note is that burst mode eliminates data bus idle time only so long as the same row is being accessed. As soon as a different row is accessed, a significant amount of idle time is introduced on the data bus 110, as shown in FIG. 3. Thus, in applications where access of the DRAM switches from row to row on a regular basis (as is often the case), there is a substantial amount of idle time on the data bus 110, even if burst mode is implemented.\nTo further improve data bus utilization, burst mode may be implemented in conjunction with a multi-bank DRAM to achieve full data bus utilization. In a multi-bank DRAM, the DRAM is divided into multiple “banks”, which may be viewed as “virtual memories” within the DRAM. Each bank may be accessed individually, and each bank has its own set of sense amps. However, all banks share the same data bus. A block diagram of a sample multi-bank DRAM 400 is shown in FIG. 4. While FIG. 4 shows a DRAM 400 having four banks 404, it should be noted that more or fewer banks may be implemented if so desired. The basic concept behind a multi-bank DRAM 400 is that higher bus utilization may be achieved by “interleaving” or alternating memory requests between the different banks 404. By interleaving the memory requests, it is possible to initiate memory access to one bank (e.g. 404(1)) while another bank (e.g. 404(2)) is busy delivering data onto the data bus. By doing so, the data bus idle time of one bank is used advantageously by the other bank to put data onto the data bus 410. Because all banks 404 are using the same data bus 410, interleaving the memory requests in this way makes it possible to keep the data bus 410 constantly filled, even when different rows are being accessed.\nTo illustrate how burst mode and interleaving may be used to achieve full data bus utilization, reference will now be made to the timing diagram of FIG. 5. In FIG. 5, a RAS1 signal is used to indicate a RAS signal applied to bank 1 404(1), while a RAS2 signal is used to indicate a RAS signal applied to bank 2 404(2), and so on. Likewise, a CAS1 signal indicates a CAS signal being applied to bank 1 404(1), while a CAS2 signal indicates a CAS signal being applied to bank 2 404(2), and so on.\nAs shown in FIG. 5, a read operation from bank 1 404(1) of the DRAM 400 is initiated by first sending an asserted RAS signal 502(1) and a row address 508(1) to bank 1 404(1). Then, at a later time, an asserted CAS signal 504(1) and a column address 510(1) are sent to bank 1 404(1). A short time thereafter, data associated with the row address 508(1) and the column address 510(1) are sent onto the data bus 410 (FIG. 4) by the sense amps 408(1) of bank 1 404(1). In the timing diagram shown in FIG. 5, it is assumed that bank 1 404(1) implements a burst mode length of two. Thus, in response to the one column address 510(1), two sets of data 506(1), 506(2) are outputted onto the data bus 410 by bank 1 404(1).\nAfter the RAS signal 502(1) is sent to bank 1 404(1) but before the CAS signal 504(1) is sent to bank 1 404(1), an asserted RAS signal 502(2) and a row address 508(2) are sent to bank 2 404(2) of the DRAM 400. In addition, an asserted CAS signal 504(2) and a column address 510(2) are sent to bank 2 404(2) at a later time. In response to these signals, bank 2 404(2) outputs data associated with the row address 508(2) and the column address 510(2) onto the data bus 410 using the sense amps 408(2). As was the case with bank 1 404(1), bank 2 404(2) also implements a burst mode length of two. As a result, two sets of data 506(3), 506(4) are outputted onto the data bus 410 by bank 2 404(2) in response to the one column address 510(2). These sets of data 506(3), 506(4) immediately follow the sets of data 506(1), 506(2) outputted by bank 1; thus, there is no idle time between the data sets.\nAfter the RAS signal 502(2) is sent to bank 2 404(2) but before the CAS signal 504(2) is sent to bank 2 404(2), an asserted RAS signal 502(3) and a row address 508(3) are sent to bank 3 404(3). A short time thereafter, an asserted CAS signal 504(3) and a column address 510(3) are sent to bank 3 404(3). In response, bank 3 404(3) outputs data associated with row address 508(3) and column address 510(3) onto the data bus 410 using the sense amps 408(3). As was the case with bank 1 and bank 2 404(2), bank 3 404(3) also implements a burst mode length of two. As a result, two sets of data 506(5), 506(6) are outputted onto the data bus 410 by bank 3 404(3) in response to the one column address 510(3). These sets of data 506(5), 506(6) immediately follow the sets of data 506(3), 506(4) outputted by bank 2; thus, there is no idle time between the data sets.\nTo finish the example, after the RAS signal 502(3) is sent to bank 3 404(3) but before the CAS signal 504(3) is sent to bank 3 404(3), an asserted RAS signal 502(4) and a row address 508(4) are sent to bank 4 404(4). A short time later, an asserted CAS signal 504(4) and a column address 510(4) are sent to bank 4 404(4). In response, bank 4 404(4) outputs data associated with row address 508(4) and column address 510(4) onto the data bus 410 using the sense amps 408(4). As was the case with the other banks, bank 4 404(4) implements a burst mode length of two. As a result, two sets of data 506(7), 506(8) are outputted onto the data bus 410 by bank 4 404(4) in response to the one column address 510(4). These sets of data 506(7), 506(8) immediately follow the sets of data 506(5), 506(6) outputted by bank 3; thus, there is no idle time between the data sets.\nWhile bank 4 404(4) is being accessed, access of bank 1 404(1) may again be initiated with a RAS signal and a row address, as shown, to extract more data from that bank. This process of interleaving accesses between the various banks 404 may continue indefinitely to continually access data from the DRAM 400. By combining burst mode operation with a multi-bank DRAM 400 as shown in this example, it is possible to achieve full data bus utilization.\nIn practice, a certain number of banks are needed to achieve full data bus utilization in a particular DRAM, where the number of banks needed is determined by certain timing parameters of the DRAM. One relevant timing parameter is the minimum time required between consecutive RAS signals to the same bank. This parameter, denoted herein as Trc, is often referred to as the RAS cycle time. Another relevant parameter is the amount of time it takes to place one base granularity of data onto the data bus. This parameter is denoted herein as Dt. To determine the number of banks needed to achieve full data bus utilization, Trc is divided by n*Dt where n is the burst length. If, for example, Trc is 80 ns and Dt is 10 ns and the DRAM is implementing a burst length of two, then the number of banks needed is 80 ns/20 ns or four. With these timing parameters (which are typical) and four banks, a DRAM can achieve full data bus utilization.\nWhile the combination of burst mode and a multi-bank DRAM 400 makes it possible to achieve 100% data bus utilization in a memory, this implementation does not come without its drawbacks. One significant drawback is that it still relies upon burst mode to achieve full data bus utilization. Because of this reliance, this implementation suffers from the same shortcoming as that experienced in regular burst mode. Namely, it increases the data granularity of the DRAM 400. Notice from the timing diagram of FIG. 5 that instead of outputting just one base granularity of data per access to each bank, the DRAM 400 outputs two (it is two in the example shown in FIG. 5 but it could more than two in other implementations). This increase in data granularity can lead to inefficiency.\nAs noted previously, in some applications, an external component (such as a CPU) may wish to access only one base granularity of data at a time. In such applications, any data provided after the first base granularity will be dropped. If burst mode is implemented in such an application, at least half of the data provided by the DRAM 400 will be dropped, which means that at most 50% efficiency can be achieved. Thus, even though burst mode combined with a multi-bank DRAM 400 may achieve 100% data bus utilization in such an application, it does not improve overall system efficiency because the extra data is not used. Consequently, the burst mode/multi-bank DRAM combination does not provide a complete solution for all possible applications.\nAs an alternative to the burst mode/multi-bank DRAM combination, a plurality of separate DRAM's may be implemented to achieve full data bus utilization. By interleaving memory requests between separate DRAM's instead of between separate banks within a single DRAM, it is possible to achieve full data bus utilization without requiring an increase in data granularity. This result comes with extra cost and complexity, however. To illustrate how separate DRAM's may be used to achieve full data bus utilization, reference will be made to FIGS. 6 and 7. Specifically, FIG. 6 shows a block diagram of a sample multi-DRAM implementation, while FIG. 7 shows a timing diagram for several read cycles of the implementation of FIG. 6.\nAs shown in FIG. 6, the sample implementation comprises a plurality of separate DRAM's 604(1), 604(2) (more than two may be implemented if so desired), with each DRAM 604 having its own separate command lines 608(1), 608(2), and address lines 606(1), 606(2). Both DRAM's share the same data bus 620, and both have 2 banks. In addition to the DRAM's 604(1), 604(2), the implementation further comprises a controller 602 for controlling the interleaving of requests between the DRAM's 604(1), 604(2). It is the responsibility of this controller 602 to manage the interleaving of memory requests such that: (1) the data bus 620 is used as fully as possible; and (2) there is no bus contention on the data bus 620.\nTo illustrate how the implementation of FIG. 6 can be used to achieve full data bus utilization, reference will be made to the timing diagram of FIG. 7. For clarity purposes, RAS 1,1 is used in FIG. 7 to indicate a RAS command applied to DRAM 1 604(1), bank 1, while RAS2,1 is used to indicate a RAS command applied to DRAM 2 604(2), bank 1, and so on. Likewise, CAS1,1 is used to indicate a CAS command being applied to DRAM 1 604(1), bank 1, while CAS2,1 is used to indicate a CAS command being applied to DRAM 2 604(2), bank 1.\nAs shown in FIG. 7, to extract data from the DRAM's 604(1), 604(2), the controller 602 first initiates a read operation on DRAM 1 604(1), bank 1. This is carried out by sending a RAS command 702(1) and a row address 706(1) to DRAM 1 604(1), bank 1. Then, a CAS command 704(1) and a column address 706(2) are sent to DRAM 1 604(1), bank 1. A short time thereafter, data 720(1) associated with the row address 706(1) and the column address 706(2) are outputted onto the data bus 620 (FIG. 6) by DRAM 1 604(1), bank 1. While the CAS command 704(1) and the column address 706(2) are being sent to DRAM 1 604(1), bank 1, the controller 602 also sends a RAS command 712(1) and a row address 716(1) to DRAM 2 604(2), bank 1. Thereafter, a CAS command 714(1) and a column address 716(2) are sent to DRAM 2 604(2), bank 1. In response to these signals, DRAM 2 604(2), bank 1, outputs data 720(2) associated with the row address 716(1) and the column address 716(2) onto the data bus 620. This set of data 720(2) immediately follows the set of data 720(1) outputted by DRAM 1 604(1), bank 1; thus, there is no idle time between the data sets.\nWhile the CAS signal 714(1) and the column address 716(2) are being sent to DRAM 2 604(2), bank 1, a RAS command 702(2) and row address 706(3) are sent to DRAM 1 604(1), bank 2, to initiate another read cycle. Thereafter, a CAS command 704(2) and a column address 706(4) are sent to DRAM 1 604(1), bank 2. In response, DRAM 1 604(1), bank 2, outputs data 720(3) associated with the row address 706(3) and the column address 706(4) onto the data bus 620. This set of data 720(3) immediately follows the set of data 720(2) outputted by DRAM 2 604(2), bank 1; thus, there is again no idle time between the data sets. To finish the example, while the CAS command 704(2) and the column address 706(4) are being sent to DRAM 1 604(1), bank 2, the controller 602 initiates a read cycle on DRAM 2 604(2), bank 2, by sending a RAS command 712(2) and a row address 716(2) to DRAM 2 604(2), bank 2. Thereafter, a CAS command 714(2) and a column address 716(4) are sent to DRAM 2 604(2), bank 2. In response to these signals, DRAM 2 604(2), bank 2, outputs data 720(4) associated with the row address 716(3) and the column address 716(4) onto the data bus 620. This set of data 720(4) immediately follows the set of data 720(3) outputted by DRAM 1 604(1), bank 2; hence, there is no idle time between the data sets. Additional read operations may be carried out in this manner to continue extracting data from the DRAM's 604(1), 604(2). As this example illustrates, by interleaving memory requests between multiple banks of multiple DRAM's, it is possible to achieve full data bus utilization, and because no burst mode is implemented, data granularity is not increased.\nWhile the multi-DRAM implementation is able to achieve full data bus utilization without increasing data granularity, it does so at a significant cost. First, due to the tight timing constraints, the DRAM's shown in FIG. 6 are very fast and very expensive. With multiple DRAM's being required to implement the system of FIG. 6, the cost of the memory system can be prohibitive. Compared to single DRAM memory systems, the cost of this multi-DRAM system can be several-fold. Also, the multi-DRAM implementation is limited in its application. By its very nature, it can be implemented only in a multi-DRAM environment. In the many applications in which it is desirable to implement just one DRAM, the multi-DRAM implementation cannot be used. In addition, this implementation can add substantial complexity to the memory access process. Because the controller 602 must concurrently control multiple DRAM's, the memory access process is much more complex and difficult to manage than in a single DRAM implementation. Overall, there is a significant price to pay for the functionality provided by the multi-DRAM configuration, and in many implementations, this price is prohibitive. Hence, the multiple DRAM approach does not provide a viable solution for all applications."}
{"text": "The present invention relates to a reactor and methods for performing single and in-situ multiple integrated circuit processing steps, including thermal CVD, plasma-enhanced chemical vapor deposition (PECVD), reactor self-cleaning, film etchback, and modification of profile or other film property by sputtering. The present invention also relates to a process for forming conformal, planar dielectric layers on integrated circuit wafers and to an in-situ multi-step process for forming conformal, planar dielectric layers that are suitable for use as interlevel dielectrics for multi-layer metallization interconnects."}
{"text": "The present invention concerns a process and a catalyst for the hydrogenation of the carbon oxides carbon monoxide and carbon dioxide. More specifically, this invention relates to a process for reacting or removing carbon oxides by hydrogenation using a catalyst comprising a support and an active component being of the composition of alloys of Ni, Fe and Co metals in oxide or reduced form.\nThe presence of carbon monoxide in feed gases is undesired in a number of processes.\nIn a fuel cell, such as a Polymer Electrolyte Membrane Fuel Cell, the presence of carbon monoxide is very critical since it can poison the noble metal electrodes used in fuel cells and therefore reduce their effectiveness.\nPreferably, the CO concentration for fuel cell feed should be less than 100 ppm, more preferably less then 50 ppm. However, the initial concentration of CO, as received from the fuel processor, can exceed 1 wt %. Therefore, further reduction of CO concentration is required. Some of the typical methods for reducing CO concentration are selective catalytic oxidation of CO, pressure swing adsorption, hydrogen separation by membrane and methanation of CO.\nSimilarly, in ammonia production plants, the presence of CO is also highly undesirable in the ammonia synthesis reactor and the carbon oxide concentrations should usually be reduced to values as low as 5-10 ppmv.\nReacting carbon monoxide and carbon dioxide with hydrogen may also be used in the preparation of methane/synthetic natural gas (SNG). SNG can be produced by gasification of biomass or coal and subsequent methanation.\nMethanation is a process where carbon oxides are reacted with hydrogen in the presence of a catalyst to produce methane and possibly smaller amounts of other lower hydrocarbons and water. In the known methanation processes, precious metals supported on a carrier as Al2O3, SiO2 or TiO2 have been used as a catalyst for CO methanation (U.S. Pat. No. 3,615,164 and WO 01/64337). These catalysts are usually able to reduce CO concentrations to values of about 500-800 ppm. Nickel-alumina catalysts, which are presently used in conventional methanation processes, contain large proportions of nickel, generally in excess of about 20 wt %. This requirement places certain limitations on practical methods for manufacturing such catalysts. In ammonia plants, methanation is mainly performed in the temperature range of 250° C. to 350° C. in the presence of Ni/Al2O3 catalysts. Modifications of such catalysts are desirable to achieve reduction and reaction at relatively lower temperatures than in the conventional design.\nAlloying an active metal with a second active or inactive metal can change the catalytic performance drastically. For example with a monometallic iron catalyst and bimetallic copper-iron catalysts (E. Boellard, F. Th. van Scheur, A. M. van der Kraan, J. W. Geus, Appl Catal. A 171 (1998) 333. It has further been demonstrated that the reduction profile, carbon monoxide chemisorption properties and the Fischer-Tropsch activity were substantially altered by adding copper to the iron phase.\nFrench patent application FR 863473-A describes a process for the production of hydrocarbons by hydrogenation of carbon monoxide in which the catalyst contains iron and nickel with an atomic ratio of 1:1.\nPatent application US 2005/0096211 describes a catalyst for fuel cell applications, which is highly selective in CO methanation, prevents the conversion of CO2 into CO and suppresses CO2 methanation. The catalyst comprises a metal selected from the group consisting of ruthenium, nickel, iron, cobalt, lead, tin, silver, iridium, gold, copper, manganese, zinc, zirconium, molybdenum, other metals that form metal-carbonyl species and combinations thereof on a support."}
{"text": "1. Field of the Invention\nThe invention relates to a multi-zone diffuser for an axial-flow turbomachine\nin which the kink angles of the diffuser inlet both at the hub and at the cylinder of the turbomachine are fixed exclusively for the purpose of evening out the total pressure profile over the duct height at the outlet from the last rotor blade row,\nin which means, in the form of streamlined struts, are provided within the deceleration zone of the diffuser for the removal of swirl from the swirling flow,\nand in which flow guide rings subdivide the diffuser in multi-duct fashion.\n2. Discussion of Background\nSuch multi-zone diffusers for turbomachines are known from EP-A 265 633. In order to meet the requirements there for the best possible pressure recovery and swirl-free diffuser outlet flow at part load, a straightening grid is provided within the diffuser and this grid extends over the complete height of the flow duct. These means for the removal of swirl involve cylindrical streamlined struts with thick straight profiles arranged uniformly around the periphery. These profiles are designed according to the knowledge available for the construction of turbomachines and are intended to be as insensitive as possible to oblique incident flow. The leading edges of these struts subjected to the incident flow are located relatively far behind the trailing edge of the last rotor blades in order to avoid excitation of the last blade row caused by the pressure field of the struts. This distance is dimensioned in such a way that the leading edge of the struts is located in a plane in which a diffuser area ratio of preferably three is present. This first diffusion zone between the blading and the streamlined struts is therefore intended to remain undisturbed because of the total rotational symmetry. The fact that no interference effects are to be expected between the struts and the blading may be attributed to the fact that the struts only become effective in a plane in which there is already a relatively low velocity level.\nBecause the opening angle of conventional highly-loaded blading of turbines far exceeds that of a good diffuser, the known diffuser is subdivided into a plurality of partial diffusers by means of flow guide rings in order to support the flow in the radial direction. These guide rings extend from a plane directly at outlet from the blading to a plane at which a diffusion ratio of three is reached, i.e. over the whole of the first diffusion zone. For vibration reasons, these guide rings should preferably be configured in one piece. This leads to a solution without a split plane, which is disadvantageous for assembly reasons. In addition, the guide rings lead to large diameters so that transport problems can arise.\nA second diffusion zone extends from the leading edge of the thick streamlined struts to the maximum profile thickness of the struts. The de-swirling of the flow is intended to take place to a major extent in this second zone and, in fact, substantially without deceleration. In a third, subsequent diffusion zone in the form of a straight diffuser, the flow--which at this time is practically swirl-free--is further decelerated.\nIn addition to maximum pressure recovery, particularly at part load, all these measures are also intended to achieve a reduction in the design length of the installation.\nIn conventional gas turbines, the flow onto the diffuser at idle has a velocity ratio c.sub.t /c.sub.n of approximately 1.2, where c.sub.t is the tangential velocity and c.sub.n is the axial velocity of the medium. This oblique incident flow leads to a reduction in the pressure recovery C.sub.p.\nIn other types of machines, such as steam turbines or gas turbines for fluidized bed firing, it is quite possible for the volume flow to be reduced down to 40% so that c.sub.t /c.sub.n ratios of up to 3 occur. In such types of machine, fixed diffuser geometry is not a possibility because the pressure recovery could even be negative. This applies even in the case where the pitch/chord ratio of the streamlined struts is 0.5. Streamlined struts with pitch/chord ratios of approximately 1, which would provide a somewhat larger pressure recovery at full load, i.e. c.sub.t /c.sub.n =approximately 0, cannot be used at all in such machines.\nThe large drop in pressure recovery may be attributed to the fact that a strong vortex forms between the outlet rotor blades and the streamlined struts in the case of the extreme relationships quoted. The vortex is bounded by the streamlined struts at which the tangential component of the velocity is dissipated. If solid particles (for example in gas turbines) or water droplets (for example in steam turbines) are entrained in the resulting reverse flow, there is an acute danger of root erosion on the blades of the last rotor row.\nA known remedy in a turbomachine of the axial type, from EP 0 417 433 A1, is to arrange at least one row of variable guide vanes in the diffuser between the means for swirl removal and the outlet rotor blades. The means for removing the swirl within the diffuser are, in this case also, streamlined struts arranged evenly around the periphery with a straight camber line and symmetrical profile and with a pitch/chord ratio between 0.5 and 1 in the center section of the flow duct. These streamlined struts extend conically in the radial direction. The intention is that the part-load behavior of the machine should be further improved by these measures for designing the diffusion."}
{"text": "Reference is made to FIG. 1 which shows a schematic diagram of a standard six transistor (6T) static random access memory (SRAM) cell 10. The cell 10 includes two cross-coupled CMOS inverters 12 and 14, each inverter including a series connected p-channel and n-channel transistor pair. The inputs and outputs of the inverters 12 and 14 are coupled to form a latch circuit having a true node 16 and a complement node 18. The cell 10 further includes two transfer (pass gate) transistors 20 and 22 whose gate terminals are coupled with a wordline node and are controlled by the signal present at the wordline node (WL). Transistor 20 is source-drain connected between the true node 16 and a node associated with a true bitline (BLT). Transistor 22 is source-drain connected between the complement node 18 and a node associated with a complement bitline (BLC). The source terminals of the p-channel transistors in each inverter 12 and 14 are coupled to receive a high supply voltage (for example, VDD) at a high voltage node VH, while the source terminals of the n-channel transistors in each inverter 12 and 14 are coupled to receive a low supply voltage (for example, GND) at a low voltage node VL. The high voltage VDD at the node VH and the low voltage GND at the node VL comprise the power supply set of voltages for the cell 10.\nIn an integrated circuit including the SRAM cell 10, this power supply set of voltages may be received at pins of the integrated circuit, or may instead be generated on chip by a voltage regulator circuit which receives some other set of voltages received from the pins of the chip. The power supply set of voltages at the nodes VH and VL are conventionally applied to the SRAM cell 10 at all times that the cell/integrated circuit is operational. It will be recognized that separate low voltage values at node VL may be provided for the sources of the n-channel MOS transistors in the inverters 12 and 14 while separate high voltage values at node VH may be provided for the sources of the p-channel MOS transistors in the inverters 12 and 14.\nThe reference above to a six transistor SRAM cell 10 of FIG. 1 for use as the data storage element is made by way of example only, it being understood to those skilled in the art that the cell 10 could alternatively comprise a different data storage element. The use of the term SRAM cell 10 will accordingly be understood to refer any suitable memory cell or data storage element, with the circuitry, functionality and operations presented herein in the exemplary context of a six transistor SRAM cell.\nReference is now made to FIG. 2 which shows a block diagram of a self-timed memory 30, for example of the static random access memory (SRAM) type using SRAM cells 10, with “w” words and “b” bits organized as a column mux of “m”. Those skilled in the art understand that self-timed memories need to support a high dynamic operating voltage range. In other words, these memories need to be functional over a wide range of supply voltages, starting from a very high operating voltage and down to a very low operating voltage. In most cases, in the low operating voltage range, it is considered acceptable if the memory achieves a lower performance (i.e., it is slower). In nominal operating voltage range, the memory needs to support a higher performance (i.e., it needs to be faster).\nThe memory 30 includes a first section 32 comprising a plurality of memory (such as SRAM) cells 10 arranged in a matrix format and which function to store data. The first section 32 includes “b” sub-sections 34 corresponding to the “b” bits per word stored by the memory. The first section 32 is arranged to store “w” words and is organized as a column mux of “m”. Thus, each of the “b” sub-sections 34 is organized in “w/m” rows with “m” columns in each row. In the first section 32, all cells 10 in a same row share a common wordline (WL) coupled to an output of a row decoder circuit 60 (well known to those skilled in the art), and all cells 10 in a same column share a common true bitline (BLT) and a common complement bitline (BLC) coupled to column circuitry 62 (which includes bitline precharge and equalization circuitry, column mux circuitry, write driver circuitry, column address decoder circuitry and input/output circuitry, each of which is well known to those skilled in the art).\nTo read data from the first section 32, the wordline of the row selected according to the row address is driven high by the row decoder circuitry 60, a column is selected in each sub-section 34 by the column address decoder and column mux in the column circuitry 62 based on the column address to connect the selected column's true bitline and complement bitline to the input/output circuitry (which, for example, will typically utilize sense amplifiers), and both the true bitline and complement bitline of the selected column in each sub-section 34 are made floating by the precharge and equalization logic in the column circuitry 62. One of the true bitline and complement bitline discharges in each sub-section 34 depending on whether the internal true node 16 or complement node 18 of the memory cell 10 in the selected row and column in the sub-section 34 stores a logic zero. This discharge is sensed by the sense amplifiers corresponding to respective sub-section 34 in the input/output logic of the column circuitry 62 to transfer the stored data state to the respective memory output pin.\nThe memory 30 includes a second section 46 including plurality of memory cells 10 arranged in a matrix format, but these cells do not function to store data. Indeed, these cells are only required, if desired, in order to have a regular layout of the memory array. The wordline ports of the memory cells 10 both the rows in this section are connected to the ground reference voltage (GND).\nThe memory 30 includes a third section 36 including plurality of memory cells 10 arranged in a matrix format, and these cells also do not function to store data. Rather, these cells in the third section 36 are used to emulate the same load on a reference wordline (REFWL), which is coupled to a reference row decoder 64 within the section 36, as is present on the actual wordlines (WL) of the first section 32. In other words, the purpose of section 36 is to emulate a total load of “b*m” columns of memory cells 10 on the reference wordline REFWL. It will be noted that the REFWL signal generated by the reference row decoder circuit 64 passes through the second section 46 to the third section 36 without being connected to cells 10 included in section 46.\nThe section 36 includes “b” sub-sections 38. Each sub-section 38 includes two rows of “m” memory cells 10. All memory cells 10 within the third section 36 either have their true bitlines and complement bitlines connected to a power supply voltage (for example, at node VH) or have them floating. The wordline ports of the memory cells 10 within one of the two rows of the first half of the total “b” sub-sections 38 (i.e., of the first “b/2” sub-sections 38) are coupled to the reference wordline signal generated by the reference row decoder circuit 64 and arriving in section 36 after having passed through the second section 46. This is done to emulate the same propagation delay corresponding to “b*m/2” columns on REFWL as is present on all the WL signals in propagating from row decoder 60 to the middle of section 32. Further, the REFWL signal which has thus reached at or about the center of the section 36 is twisted back and returned towards reference row decoder circuit 64. This returning REFWL signal is connected to the other of the two rows of the first half of the total “b” sub-sections 38 (i.e., of the first “b/2” sub-sections 38), eventually reaching the second section 46 again after experiencing a propagation delay corresponding to travelling across “b*m” columns—same as that experienced by the signal WL in propagating from row decoder 60 to the column farthest from the row decoder 60 at the end of section 32. The reference wordline of the memory cells 10 in both rows within other “b/2” sub-sections 38 (i.e., later “b/2” sub-sections 38) is coupled to a ground supply voltage (at the node VL) because these sub-sections 38 are present in the memory only for maintaining regularity and rectangular shape of the array of the memory cells 10, and so the memory cells 10 in these sub-sections 38 are deactivated permanently by connecting their wordline ports to a ground supply voltage (for example, at the node VL).\nThe memory 30 further includes a fourth section 40 including a plurality of read timer cells 42 and load cells 44 arranged in a matrix format: “w/m” rows and one column. The timer cells 42 and load cells 44 each have a configuration similar to a memory cell 10, but with different connections as shown, in an exemplary manner, in FIGS. 3A-3C for timer cells 42 and FIGS. 3D-3E for load cells 44.\nThe timer cells 42 are essentially pre programmed memory (for example, SRAM) cell like elements that are built from the same devices as used by the memory cells 10 in section 32. These timer cells 42 operate to discharge the reference true bitline (REFBLT) in response to arrival of a reference wordline (REFWL) signal (with their discharge time being indicative of time required to discharge of actual bitline until safe detection level during a read operation). It may also be noted that the timer cell configurations in FIGS. 3A and 3B discharge both reference true bitline (REFBLT) and reference complement bitline (REFBLC), and are helpful in using double the number of discharge cells with double the load on the reference bitline for improved statistical variability without any area impact. However, going further, for the purpose of description of functionality only, a description is provided for the use of discharge of only the reference true bitline (REFBLT) at the next logic stage, without any loss in the description of functionality.\nThe load cells 44 are elements similar to timer cells 42, with the difference that either their reference wordline (REFWL) ports only or both the reference wordline (REFWL) ports and the gate terminal of the pulldown transistor in series with the passgate transistor coupled with reference true bitline (REFBLT) are grounded, so that they do not discharge REFBLT and REFBLC and only serve to match the load of actual bitlines (BLT/BLC) on REFBLT and REFBLC. The wordlines WL generated in the row decoder circuitry 60 simply pass through this section 40 in order to reach the first section 32.\nWith respect to the timer cells 42, the reference true bitlines (REFBLT) and reference complement bitlines (REFBLC) are shared with each other, respectively, along a column and the reference wordlines (REFWL) for the timer cells 42 are shared with each other and coupled to the reference wordline (REFWL) node at the point emerging from the second section 46 after having traversed through and returned after half the columns in the memory in third section 36. With respect to the load cells 44, the reference wordline (REFWL) ports are grounded as already described earlier. Furthermore, with respect to the timer cells 42 and load cells 44, the reference true bitlines (REFBLT) and reference complement bitlines (REFBLC) are shared with each other, respectively, along a column and connected to the column circuitry 62.\nThere are a total of “w/m” timer cells 42 and load cells 44, in order to emulate same load on the reference true and complement bitlines within section 40 as is present on the true and complement bitlines within first section 32. A certain number of these “w/m” cells are timer cells 42, and the remaining are load cells 44. Multiple timer cells 42 are used in order to discharge the reference true bitline (REFBLT) at a much faster rate as compared to the rate at which the selected actual memory cells 10 discharge the actual bitlines (BLT or BLC depending on data stored in the respective memory cells 10) during a read operation. The rate of discharge for reference true bitline (REFBLT) is desired to be much faster than the rate of discharge of the actual bitline (true or complement) so that the node REFBLT is able to discharge to a level detectable by an inverter circuit contained within the column circuitry 62 (for example, almost VDD/2) in the same time in which a statistically worst memory cell 10 is able to discharge the actual bitline (BLT or BLC) to a voltage level beyond the offset of the sense amplifiers (contained in column circuitry 62). Also, the multiple timer cells 42 help in reducing statistical variability of the access and cycle time of read operation. This detection of REFBLT discharge is propagated by subsequent logic to generate an enable signal which activates the sense amplifiers (also contained in the column circuitry 62) during a read operation. Thus, the intention of this operation is to time the activation of the sense amplifier at an optimum time permitting a certain statistically worst memory cell 10 in section 32 to clear the sense amplifier offset for reading from that memory cell 10.\nA more detailed description of memory operation is now provided. Before any read cycle begins, all memory bitlines and the reference bitlines are precharged to logic high (VDD), and all memory wordlines (WL) and the reference wordline (REFWL) are driven to logic low (GND). At the start of a valid read operation characterized by the “clock” edge when the “chip select” signal is asserted for enabling the memory and the “write enable” signal is asserted to the level corresponding to the read operation, a clock generator triggers the internal clock signal at the arrival of the “clock” edge (either rising or falling edge depending on the functionality of the memory). The internal clock signal causes the following operations to be performed more or less concurrently: a) driving of one of the “w/m” wordlines (WL), selected depending on row address, to logic high (VDD); b) driving of the reference wordline (REFWL) to logic high (VDD); c) turning off of precharge of the reference bitlines (REFBLT, REFBLC) and of the bitlines (BLT, BLC) of one of the “m” columns of each of “b” bits, selected depending on column address; d) connecting one of the “m” true bitline (BLT) and complement bitline (BLC) pairs in each of the “b” bits to the inputs of the sense amplifier circuitry (contained within the column circuitry 62).\nThe triggering of above operations in turn starts the following operations (performed more or less concurrently): a) the rising of the selected wordline starts an operation in each selected memory cell 10 (in the first section 32) to discharge either the true bitline (BLT) or complement bitline (BLC) based on the stored data; and b) the rising of reference wordline (REFWL) starts an operation in the timer cells and load cells of third section 40 to discharge the reference true bitline (REFBLT).\nIt will be noted that in each bit (I/O), there is only a single memory cell 10 which causes a discharge of one of the true and complement bitlines (BLT, BLC), but there are multiple timer cells 42 in a column causing a discharge of the reference true bitline (REFBLT). Thus, over a same time period in which one of the true and complement bitlines (BLT, BLC) in any bit (I/O) discharges by a first voltage amount (e.g., x mV), the reference true bitline (REFBLT) would be expected to discharge by a second voltage (e.g., y mV), where y>x, because the number of timer cells 42 discharging the reference true bitline (REFBLT) is greater than the single actual memory cell 10 discharging the bitlines (BLT or BLC) in any bit (I/O). It can also be observed that the time it takes for bitline discharge (i.e., data to be read from any memory cell 10 of the first section 32) is much more variable than what it takes for true reference bitline (REFBLT) discharge from the timer cells 42 of the fourth section 40 because of the presence of multiple timer cells (the standard deviation of the read current of “n” memory cells is known to be “sigma/squareroot(n)” where “sigma” is the standard deviation of a single memory cell 10).\nIn the memory of FIG. 2, as per the prior art, both the wordline WL and reference wordline REFWL are driven by similarly sized drivers, to a full logic high (VDD), causing the discharge of one of the bitlines (either BLT or BLC depending on data present in the selected memory cell 10 in any bit (I/O)) as well as the true reference bitline REFBLT. The discharge of reference bitline further generates the sense amplifier enable signal. The generated sense amplifier enable signal activates the sense amplifier for each bit (I/O) so as to allow the sense amplifier circuitry to resolve the discharge of one of the true or complement bitlines (BLT, BLC) of the selected column in that bit (I/O) to a logic low (GND) or logic high (VDD). The data resolved by the sense amplifiers of different bits (I/Os) is transferred to memory output (“Q”) of respective bits, thus completing the read operation.\nIt is desirable to have the timer cells 42 designed and their number chosen such that, in about the same time that a statistically worst memory cell 10 takes to discharge one of the true and complement bitlines (BLT, BLC) beyond the sense amplifier offset voltage, on a certain worst process (P), voltage (V) and temperature (T) condition, the multiple timer cells 42 are able to discharge the true reference bitline REFBLT to a level detectable by a simple detector circuit such as an inverter within the column circuitry 62. That way, the discharge of true reference bitline REFBLT can be detected by the column circuitry 62 to generate the sense amplifier enable signal at an optimum time for performing the read operation successfully and with a least read access and cycle time.\nReference is now made to FIG. 4 which presents a timing diagram illustrating the read operation. From FIG. 4, it can be observed that in order to design a robust memory (i.e., one capable of successful read/write operations in spite of high statistical variability in electrical characteristics of the memory cell 10), it is important to tune the delay duration “TREFBLT” (measuring the delay from initiation of the discharge of the true or complement bitline (BLT, BLC) of any bit (I/O) to a full voltage swing discharge of the true reference bitline REFBLT) in such a way that a read from a statistically worst memory cell 10 is able to provide sufficient voltage differential (partial voltage swing) between the true and complement bitlines (BLT, BLC) to clear the sense amplifier voltage offset for correct operation of sense amplifier resulting in sensing of the appropriate read data. Thus, if “iread” is the read current of a statistically worst memory cell 10, and “irefread” is the current of the “n” timer cells 42, the ideal scenario is to have the ratio “irefread/iread” substantially constant across the process, voltage and temperature conditions of the memory. The ratio should be such that for any process, voltage and temperature condition, the time it takes for one of the true and complement bitlines (BLT, BLC) to discharge beyond the sense amplifier offset is equal to the time it takes for the true reference bitline REFBLT to discharge almost fully and be detected by a simple detector circuit such as an inverter which is used to generate the sense amplifier enable (SAEN) signal.\nThe issue with the prior art design described above is that the read current (i.e., iread) for a statistically worst memory cell 10 in section 32 is produced by the pass-gate and pulldown transistor devices (see, FIG. 1 and the n-channel passgates 20 and 22 and the n-channel pulldowns within latches 12 and 14) whose threshold voltages (“vtpg” and “vtpd”, respectively) are much higher than those of any nominal memory cell 10 (i.e., memory cell with average read current) within section 32. However, the current which discharges the true reference bitline REFBLT is produced by multiple timer cells 42. Those multiple timer cells 42 coupled together in parallel in section 40, while reducing the statistical variability and hence resulting in discharge current of the timer cell which more nearly matches that of the passgate and pulldown devices for memory cells 10 in section 32 having nearly nominal threshold voltages with multiplicity of number of number of timer cells, is not representative of the statistically worst memory cell 10 (in terms of voltage characteristics). This condition may lead to non-optimal voltage scaling characteristics of sense amplifier enable time and hence also of memory read access time.\nThe reason for the difference in voltage characteristics mentioned above is that the current of a mosfet as a function of operating voltage is such that the change in current is much greater with change in voltage when the operating voltage of the memory is nearer to the transistor threshold voltage, as compared to when operating voltage of the memory is much higher than the transistor threshold voltage. Thus, those skilled in the art will recognize that as per the design described above, if “n1” timer cells 42 are required for qualifying the statistically worst memory cell 10 read current at a certain low operating voltage (i.e., to achieve a desired irefread/iread ratio), then at a higher operating voltage, because of the discharge current of timer cells 42 increasing by lesser percentage compared to the discharge current of statistically worst memory cell 10, the same “n1” timer cells would lead to a lower “irefread/iread” ratio (i.e., a successful but slower operation leading to sub-optimal access time of read operation). This means that as per design approach described above (primary design parameter is the electrical characteristics of read timer cell 42 and its number), if a memory read self timing is designed to be functional at a certain low voltage using a certain number of timer cells (say “n1”), that same number “n1” timer cells will lead to a much slower operation and hence higher read access and cycle time what is possible to achieve when the memory is operated at a higher voltage. In other words, the performance (read access time and cycle time) will be sub-optimal at higher operating voltages. The amount of extra margin that is undesirably introduced at those higher operating voltages increases as the designed-to minimum operating voltage is lowered.\nTo address this issue, one prior art circuit uses two different numbers of timer cells 42 (say “n1” and “n2”), to be selectively chosen based on operating voltage of the memory. For example, “n1” timer cells 42 would be used over one range of operating voltage and “n2” timer cells 42 would be used over another range of operating voltages. To assist the operation for selecting a number of timer cells 42 for use, the memory device would include a low voltage control signal input (for example, LV) for selecting the number of timer cells 42. If the control signal input is in a first logic state, then n1 timer cells 42 are enabled for use in discharging the true reference bitline REFBLT, while if the control signal input is in a second logic state, then n2 timer cells 42 are enabled for use in discharging the true reference bitline REFBLT.\nLimitations of this prior art solution include: a more complicated system design resulting from having to generate and process the low voltage control signal input (LV) based on operating voltage; and voltage scaling of memory read access and cycle time that are not seamless across the entire operating voltage range because there will be an abrupt change in read performance when the operating voltage changes across the threshold point for low voltage control signal (LV) and the low voltage control signal transitions in response thereto.\nMore generally speaking, the prior art solutions consider two options.\nIn a first option, the designer may decide the number of timer cells 42, capacitance on the reference bitline and logic delay subsequent to reference bitline when the reference wordline REFWL level is full logic high level “VDD”, based on the required “wordline on to sense amplifier enable” time for a statistically worst memory cell 10 to generate a voltage differential greater than sense amplifier offset at the inputs of the sense amplifier, at the minimum required memory functional voltage at the worst case temperature and for the worst case process corner, and operate with the same setting on other voltages.\nIn a second option, the designer may decide the number of timer cells 42, capacitance on the reference bitline and logic delay subsequent to reference bitline at multiple voltage points (typically two) when the reference wordline REFWL=VDD, after sub-dividing the required memory functional voltage range into smaller voltage ranges (typically two), tune the parameters to obtain required “wordline on to sense amplifier enable” time at the lower voltage limit of each range, and then control the selection of respective setting with control pins required to be asserted/deasserted corresponding to the voltage range of operation at any point of time.\nBoth of the above solutions have shortcomings as described earlier. A need therefore exists in the art to address the foregoing problems with respect to self-timed memory read operation over a wide range of supply voltage. Such a memory will have optimum read access and cycle time at all voltages in the supported voltage range including the nominal (higher) voltage range. Thus, such a memory will have a very low read access and cycle time value in nominal (higher) voltage range required during high frequency operations while still remaining functional for read in a lower operating voltage, without any control signal requirements from the system."}
{"text": "1. Field of the Invention\nThe present invention relates generally to information retrieval systems and, more particularly, to natural language information retrieval systems.\n2. Related Art\nAs information systems become increasingly interconnected through intranets and internets, the main problem with the search for information has shifted from determining whether the requisite information exists to determining how to locate such information.\nTo explore the Internet or some other large database, well-known browsers and search engines are available. Unfortunately, currently existing search engines generally require the use of expressions and search methods with which the user has to be familiar, and generally require the user to enter keywords assumed to be related to the information that the user seeks. Moreover, a typical search returns a vast amount of information, much of it being irrelevant to the user. The user is then required to find the few relevant documents from the search results. Furthermore, it is not uncommon to not locate interesting documents because the user did not use keywords that corresponded to the words or word forms as written in those documents. Various conventional techniques have been developed to improve the recall of these and other information retrieval systems. However, these techniques have many drawbacks that limit their effectiveness in improving the ability of information retrieval systems to identify all available information related to a desired search topic.\nA primary drawback is that the conventional method for expressing the content of a text is a single word extraction. Conventional information retrieval methods rely on word stemming or rooting, skip word filtering, and proximity measures. Typically, conventional systems stem or root the words that occur in the text and subsequently filter out all stems or roots that appear in a predetermined skip list. The skip list contains words that have little or no predictive value. Such words include function words, such as articles, pronouns, prepositions and other frequently used words. The end result of these conventional methods is a keyword list containing a list of single non-trivial words occurring in the document, optionally ordered by their frequency in the text. The keywords in the keyword list are in their stemmed or rooted form, and accompanied by their offset values or similar location markers.\nFor computer applications that require an intelligent representation of the content of a text, these conventional methods are inadequate. For example, meaningful units of content generally consist of more than one single word as provided by the conventional information retrieval systems. For example, in a keyword list containing, among others, \"Amsterdam\", \"Rotterdam\", \"Marathon\", and \"Airport\", the informational content of the text, the Rotterdam Marathon and Amsterdam Airport, is lost. Likewise, in a query \"are the Antwerp Yellow Pages in the Web yet?\", the crucial phrase evidently is \"Antwerp Yellow Pages\". This phrase needs to be parsed and processed in a way that retains the informational content of the query. Specifically, only those documents that literally match on \"Yellow Pages\" and not just any occurrence of either \"page\" and \"yellow\" separately, as well as on \"Antwerp\" in a premodifying or postmodifying position, rather than just in any location in the text, should be retrieved from the searched database. However, conventional information retrieval systems typically yield the query \"Antwerp or Yellow or Page\" which may retrieve, among others, documents on Flemisch paper factories. Furthermore, the single keywords \"yellow\" and \"page\" fail to express the notion that emerges from their combination and preservation of the plural form; that is, \"Yellow Pages\".\nIn addition, the single keyword lists used in conventional information retrieval systems do not merge expressions that are different in form but share the same reference. A method that ignores synonyms, hyponymys, name variants, frequent misspellings, and other semantic relations, fails to give a proper representation of the document's content. For example, a text dealing with the wife of the current president of the United States may contain any number of references to that person, ranging from \"the President's wife\" and \"The First Lady\" to \"Mrs. Clinton\" and \"Hilary.\" Conventional information retrieval systems ignore these synonymous expressions, failing to give a proper representation of a text's contents.\nAnother drawback to conventional systems is that their mechanical application of a skip word list ignores the content representation of the text. For example, the individual word \"page\" may be a skip word, but in combinations like \"Yellow Pages\" or even \"The Sports Pages,\" it should be preserved. Likewise, skipping \"first\" from \"First Lady\" leads to a loss of the essence of the expression.\nWhat is needed, therefore, is an apparatus and method for efficiently retrieving information that accurately represents the content of both the text being searched, and the user's query, in such a way that the two can be more effectively matched."}
{"text": "Process control relates to the control of the parameters of a process, such as a manufacturing process. The parameters can include duration, temperature, pressure, speed, quantity, and/or direction of motion, for example. In other processes, such as in those performed in the course of using and maintaining information systems, management systems, and the like, the parameters can relate to paper documents or machine-readable media, or to intangible quantities such as units of time or units of money such as a credit limit, or to a result such as can be created by a printer, a visual display module, or an aural signal. Process control systems can be used to help manage production, monitor and maintain equipment, and/or perform business functions such as maintaining inventory records, performing accounting, and tracking sales."}
{"text": "Ebola Virus\nMinus-strand (−) RNA viruses are major causes of human suffering that cause epidemics of serious human illness. In humans the diseases caused by these viruses include influenza (Orthomyxoviridae), mumps, measles, upper and lower respiratory tract disease (Paramyxoviridae), rabies (Rhabdoviridae), hemorrhagic fever (Filoviridae, Bunyaviridae and Arenaviridae), encephalitis (Bunyaviridae) and neurological illness (Bomaviridae). Virtually the entire human population is thought to be infected by many of these viruses.\nThe Ebola virus comes from the Filoviridae family, similar to the Marburg virus. It is named after the Ebola River in Zaire, Africa, near where the first outbreak was noted by Dr. Ngoy Mushola in 1976 after a significant outbreaks in both Yambuku, Zaire (now the Democratic Republic of the Congo), and Nzara, in western Sudan. Of 602 identified cases, there were 397 deaths.\nThe two strains identified in 1976 were named Ebola-Zaire (EBO-Z) and Ebola-Sudan (EBO-S). The outbreak in Sudan showed a lower fatality rate—50%—compared to the 90% mortality rate of the Zaire strain. In 1990, a second, similar virus was identified in Reston, Va. amongst monkeys imported from the Philippines, and was named Ebola-Reston.\nFurther outbreaks have occurred in Zaire/Congo (1995 and 2003), Gabon (1994, 1995 and 1996), and in Uganda (2000). A new subtype was identified from a single human case in the Côte d'Ivoire in 1994, EBO-CI.\nOf around 1500 identified human Ebola infections, two-thirds of the patients have died. The animal (or other) reservoir which sustains the virus between outbreaks has not been identified.\nAmong humans, the Ebola virus is transmitted by direct contact with infected body fluids such as blood.\nThe incubation period of Ebola hemorrhagic fever varies from two days to four weeks. Symptoms are variable too, but the onset is usually sudden and characterised by high fever, prostration, myalgia, arthralgia, abdominal pains and headache. These symptoms progress to vomiting, diarrhea, oropharyngeal lesions, conjunctivitis, organ damage (notably the kidney and liver) by co-localized necrosis, proteinuria, and bleeding both internal and external, commonly through the gastrointestinal tract. Death or recovery to convalescence occurs within six to ten days of onset of symptomology.\nThe development of a successful therapeutic for Ebola virus is a long-sought and seemingly difficult endeavor. Although they cause only a few hundred deaths worldwide each year, filoviruses are considered a significant world health threat and have many of the characteristics commonly associated with biological weapons since they can be grown in large quantities, can be fairly stable, are highly infectious as an aerosol, and are exceptionally deadly. Filoviruses are relatively simple viruses of 19 Kb genomes and consist of seven genes which encode nucleoprotein (NP), glycoprotein (GP), four smaller viral proteins (VP24, VP30, VP35 and VP40), and the RNA-dependent RNA polymerase (L protein) all in a single strand of negative-sensed RNA. Administration of type I interferons, therapeutic vaccines, immune globulins, ribavirin, and other nucleoside analogues have been somewhat successful in rodent Ebola virus models, but not in nonhuman primate infection models.\nIn view of the severity of the diseases caused by (−) RNA viruses, in particular members of the Filoviridae family of viruses, and the lack of effective prevention or therapies, it is therefore an object of the present invention to provide therapeutic compounds and methods for treating a host infected with a (−) RNA virus.\nsiRNA\nDouble-stranded RNA molecules (dsRNA) have been shown to block gene expression in a highly conserved regulatory mechanism known as RNA interference (RNAi). WO 99/32619 (Fire et al.) discloses the use of a dsRNA of at least 25 nucleotides in length to inhibit the expression of genes in C. elegans. dsRNA has also been shown to degrade target RNA in other organisms, including plants (see, e.g., WO 99/53050, Waterhouse et al.; and WO 99/61631, Heifetz et al.), Drosophila (see, e.g., Yang, D., et al., Curr. Biol. (2000) 10:1191-1200), and mammals (see WO 00/44895, Limmer; and DE 101 00 586.5, Kreutzer et al.). This natural mechanism has now become the focus for the development of a new class of pharmaceutical agents for treating disorders that are caused by the aberrant or unwanted regulation of a gene.\nRecent reports have indicated that in vitro, RNAi may show some promising in reducing Ebola replication and providing protection in guinea pigs (Geisbert, et al., The Journal of Infectious Diseases, 193 (2006), 1650-1657). However, the RNAi agents examined were not designed against all known Ebola strains and were not selected for stability and other properties needed for in vivo therapeutic RNAi agents. Accordingly, despite significant advances in the field of RNAi, there remains a need for an agent that can selectively and efficiently silence a gene in the Ebola virus using the cell's own RNAi machinery that has both high biological activity and in vivo stability, and that can effectively inhibit replication of the Ebola virus for use in treating pathological processes mediated by Ebola infection."}
{"text": "Antibody conjugates of cytotoxic drugs are being developed as target-specific therapeutic agents. Antibodies against various cancer cell-surface antigens have been conjugated with different cytotoxic agents that inhibit various essential cellular targets such as microtubules (maytansinoids, auristatins, taxanes: U.S. Pat. Nos. 5,208,020; 5,416,064; 6,333,410; 6,441,163; 6,340,701; 6,372,738; 6,436,931; 6,596,757; 7,276,497), DNA (calicheamicin, doxorubicin, CC-1065 analogs; U.S. Pat. Nos. 5,475,092; 5,585,499; 5,846,545; 6,534,660; 6,756,397; 6,630,579). Antibody conjugates with some of these cytotoxic drugs are actively being investigated in the clinic for cancer therapy (Richart, A. D., and Tolcher, A. W., 2007, Nature Clinical Practice, 4, 245-255).\nThe antibody-cytotoxic agent conjugates typically are prepared by the initial modification of reactive moieties on antibodies, such as lysine amino groups, or cysteine groups (generated by reduction of native disulfide bonds or by engineering of additional non-native cysteine residues on to antibodies using molecular biology methods). Thus antibodies are first modified with a heterobifunctional linker reagent, such as those previously described, exemplified by SPDB, SMCC and SIAB (U.S. Pat. No. 6,913,758 and U.S. Patent Publication No. 20050169933) to incorporate a linker with a reactive group such as mixed pyridyldisulfide, maleimide or haloacetamide. The incorporated reactive linker group in the antibody is subsequently conjugated with a cytotoxic agent containing a reactive moiety such as a thiol group. Another conjugation route is by reaction of a cytotoxic agent derivative containing a thiol-reactive group (such as haloacetamide, or maleimide) with thiol groups on the cell-binding agent. Thiol groups are incorporated on cell-binding agents such as an antibody by reduction of native disulfide residues (R. Singh et al., Anal. Biochem., 2002, 304, 147-156), or reduction of incorporated disulfide moieties (via SPDP, succinimidyl 3-(2-pyridyldithio)propionate, followed by reduction with dithiothreitol, D. G. Gilliland et al., Proc. Natl. Acad. Sci. USA., 1980, 77, 4539-4543), or by incorporation of additional non-native cysteine residues (J. B. Stimmel et al., J. Biol. Chem., 2000, 275, 30445-30450), or incorporation of thiol groups by reaction with 2-iminothiolane (R. Jue et al., Biochemistry, 1978, 17, 5399-5406), or methyl 3-mercaptopropionimidate ester (T. P. King et al., Biochemistry, 1978, 17, 1499-1506).\nThe antibody-cytotoxic agent conjugates with disulfide or thioether linkages are cleaved intracellularly, presumably in lysosomes, to deliver the active cytotoxic agent inside the cancer cell (H. K. Erickson et al., 2006, Cancer Research, 66, 4626-4433). In addition to the killing of target cells, antibody-cytotoxic agent conjugates with reducible disulfide linkage also kill proximate antigen-negative cells in mixed populations of antigen-negative and antigen-positive cells in vitro and in vivo in xenograft models, suggesting the role of target-cell released cytotoxic agent in improving potency against neighboring non-antigen-expressing cells in tumors with heterogeneous antigen expression (Y. V. Kovtun et al., Cancer Research, 2006, 66, 3214-3221).\nAlthough, antibody-cytotoxic drug conjugates show cell killing activity in vitro and anti-tumor activity in vivo, their potency is diminished in many cases, especially when the antigen expression on the target cancer cell is low, or when the target cells are resistant to the treatment. This is often the case in the clinical setting, resulting in low to modest anti-tumor activity in patients. A potential approach to try to circumvent resistance is to synthesize new drugs that bear hydrophilic or lipophobic functionalities (see G. Szokacs et al., Nature Reviews, 5; 219-235, 2006). However, this process is cumbersome and several analogs have to be synthesized, and often modification in the structure of the drug results in loss of biological activity. Thus, there is a need for a different approach.\nThe method described in the art for preparing a cytotoxic conjugate of a cell binding agent and a drug via non-cleavable linker requires two reaction steps (U.S. Pat. No. 5,208,020 & US Publication No. 2005/0169933). First, the cell binding agent, such as an antibody, is modified with a bifunctional crosslinker that undergoes reaction with the reactive groups of the cell binding agent, such as the amine group on lysine residues or the sulfhydryl group on cysteine residues, to form covalent chemical bonds. Following modification of the cell binding agent, the product is purified to separate the desired modified cell binding agent from unreacted crosslinker. In a second step, known as a conjugation step, a reactive drug derivative, such as a thiol-containing maytansinoid, is added to the modified cell-binding agent for reaction with the modified cell-binding agent. Following this reaction, an additional purification is required to remove any unreacted drug species and other byproducts from the final conjugate. These multiple reaction and purification steps result in low yield of the final conjugate and can be expensive and cumbersome when one considers implementing these steps on a large scale. An additional drawback to these methods is the conjugate heterogeneity that is introduced when unreacted crosslinker remains linked to the cell-binding agent without attached drug. The unreacted crosslinker can then undergo additional side reactions such as hydrolysis and intramolecular or intermolecular reactions. There is therefore a need for functionalized, reactive drug derivatives, such as maytansinoids, that can be covalently linked via a non-cleavable bond to a cell binding agent, such as an antibody in essentially one reaction step."}
{"text": "Prior to the conception and development of the present invention, as is generally well known in the prior art, light bulbs are used in practically every room of a home or office and must be replaced every so often. Climbing up a ladder or balancing on a chair to change burned out bulbs in a high ceiling fixture can be quite difficult and dangerous for anyone, even possibly resulting in injury to those involved."}
{"text": "Many conveyor belts are known to the art which have been designed in a manner such as to enable them to be adapted for the transport of different goods. For example, there is known a conveyor-belt comprising a rubber or plastic-coated belt for transporting granular or particulate material. Such a conveyor belt is limited with respect to the angle of inclination of the belt, owing to the inherent angle of repose of said material.\nTo enable material to be transported on such a belt at an angle greater than the repose angle, it has been proposed that the belt be provided with a plurality of preferably uniformally spaced ridges or transverse ribs.\nAlso known is a conveyor belt which comprises a plurality of links where the mutually opposite end parts of respective links are pivotally connected to respective end parts of adjacent links, said links forming an endless belt which is guided by drive wheels and guide-rollers or guide-wheels. One advantage with conveyor belts of this latter kind is that they are relatively resistant to wear by the goods transported thereon; another advantage is that the links can be provided with side-boards such as to permit the volume of material per unit of time to be increased."}
{"text": "The concept of multi-analyte sensing using array based sensors (Chem. Rev. 100, 2595-2626, (2000)) has opened up a wide field of technologies in detecting and analyzing specific components (analytes) in a mixture of unknowns. Such technologies benefit industries ranging across the medical, biological, environmental, and consumer sectors. For example, the medical industry depends on analysis for the detection of metabolites, drugs, and glucose; the biological sector needs to detect amino acids, cell components; environmentalists have a need to know the level of gaseous components in water or air; while consumers may want to regularly test for levels of carbon monoxide in houses, airborne allergens, or hardness of water.\nThe basic principles of microarray assays were already described by the end of the eighties (J. Pharm Biomed Anal 7, 155-168, (1989)). Interest increased dramatically with the development of DNA chip technology. The invention and demonstration in the early 1990s (Science, 251, 767-773, (1991)) of high-density arrays formed by spatially addressable deposition of sensors on a two-dimensional solid support has enhanced and simplified the process of array based sensor technologies. The key to current microarray technology is the placement of receptors at predetermined locations on a microchip in a “spatially addressable” manner. The presence or absence of an analyte is then discerned by monitoring a specific location on a sensor array of receptors. All of these systems require preparing a sensor array with a plurality of receptors at predetermined sites that involve complex and expensive processing steps.\nRecent technologies have used various approaches to fabricate microarrays. For example, U.S. Pat. Nos. 5,412,087, and 5,489,678 demonstrate the use of a photolithographic process for making peptide and DNA microarrays. These patents teach the use of photolabile protecting groups to prepare peptide and DNA microarrays through successive cycles of deprotecting a defined spot on a 1 cm×1 cm chip by photolithography, then flooding the entire surface with an activated amino acid or DNA base. Repetition of this process allows construction of a peptide or DNA microarray with thousands of arbitrarily different peptides or oligonucleotide sequences at different spots on the array. This method is expensive. Park et al. (Science 276:1401 (1997)) have disclosed a lithographic method for producing an array of nanometer-sized holes using polystyrene-polybutadiene copolymer masks in reactive ion etching of silica nitride. This multi-step method is capable of producing arrays of picoliter-sized holes that are typically 20 nanometers in diameter and 20 nanometers deep with a spacing of 40 nanometers. Hole densities of up to 1011 holes/cm2 are disclosed. The range of sizes and spacings of the holes produced by this method is limited by the size of the copolymer microdomains. Uniformity of hole size and spacing is difficult to maintain with this method due to difficulties in controlling the etching method employed to form the holes.\nBecause the number of bioactive probes to be placed on a single chip usually runs anywhere from 1,000 to 100,000 probes, the spatially addressable method is intrinsically expensive regardless of how the chip is manufactured. An alternative approach is the use of fluorescent dye-incorporated polymeric beads to produce biological multiplexed arrays. U.S. Pat. No. 5,981,180 discloses a method of using color coded beads in conjunction with flow cytometry to perform multiplexed biological assay. Microspheres conjugated with DNA or monoclonal antibody probes on their surfaces were dyed internally with various ratios of two distinct fluorescence dyes. Hundreds of “spectrally addressable” microspheres were allowed to react with a biological sample and the “liquid array” was analyzed by sequentially passing microspheres through a flow cytometry cell to decode sample information. U.S. Pat. No. 6,023,540 discloses the use of fiber-optic bundles with pre-etched microwells at the distal ends to assemble dye loaded microspheres. The surface of each spectrally addressed microsphere was attached with a unique bioactive agent and thousands of microspheres carrying different bioactive probes combined to form “beads array” on pre-etched microwells of fiber optical bundles. More recently, a novel optically encoded microsphere approach was accomplished by using different sized zinc sulfide-capped cadmium selenide nanocrystals incorporated into microspheres (Nature Biotech. 19, 631-635, (2001)). Given the narrow spectral band width demonstrated by these nanocrystals, this approach significantly expands the spectrally addressable barcoding capacity in microspheres.\nEven though the “spectrally addressed microsphere” approach does provide an advantage in terms of its simplicity over the old fashioned “spatially addressable” approach in microarray making, there are still needs in the art to render the manufacture of microarrays less difficult and less expensive.\nU.S. patent application Ser. No. 09/942,241 discloses a microarray that is less costly and easier to prepare than those previously disclosed, because the support need not be modified, even though the microspheres remain immobilized on the substrate. The disclosed microarray includes microspheres dispersed in a fluid containing a gelling agent or a precursor to a gelling agent, wherein the microspheres are immobilized at random positions on the substrate. The substrate is free of receptacles designed to physically or chemically interact with the microspheres. Disclosed are a unique coding composition and technology to prepare a microarray on a substrate that does not require placement of microspheres at predetermined locations. Various coating methods are taught, but there is exemplified machine coating, whereby a support is layered with a fluid coating composition comprising microspheres dispersed in gelatin. Immediately after coating, the support is passed through a chill-set chamber in the coating machine where the gelatin undergoes rapid gelation and the microspheres are immobilized.\nAlthough the disclosure of this patent application provides a manufacturing advantage over other existing technologies, some limitations need to be overcome. By moving from spatially addressable to randomly positioned microspheres, the information content contained within each bead necessarily must be extracted using a new analysis technology that is not preset-positionally dependent. Furthermore, the colors and color levels need to be accessed uniquely to correlate the tag to the analyte.\nIt is also known (Nature Biotech. 19, 631-635, (2001)) that the number of different color codes in spectrally addressable microspheres for use in multi-analyte sensing follows the relationship:Number of optical codes=(nm−1), where                m=color types, and        n=color intensity levels        \nFor example, 2 colors, with 4 intensity levels each should result in 42−1=15 codes. Hence, in order to sense a large number of analytes using several color types and numerous color levels, there exists a need for analysis methods to differentiate small changes in color types and color levels, on a micrometer scale."}
{"text": "Field of the Invention\nThe present invention relates to devices and methods for energetically coupling a detonator cord to one or more shaped charges.\nDescription of the Related Art\nHydrocarbons, such as oil and gas, are produced from cased wellbores intersecting one or more hydrocarbon reservoirs in a formation. These hydrocarbons flow into the wellbore through perforations in the cased wellbore. Perforations are usually made using a perforating gun loaded with shaped charges. The gun is lowered into the wellbore on electric wireline, slickline, tubing, coiled tubing, or other conveyance device until it is adjacent to the hydrocarbon producing formation. Thereafter, a surface signal actuates a firing head associated with the perforating gun, which then detonates the shaped charges. Projectiles or jets formed by the explosion of the shaped charges penetrate the casing to thereby allow formation fluids to flow through the perforations and into a production string.\nReferring now to FIG. 1A, there is shown a conventional perforating gun 10. The gun 10 includes a charge strip or tube 12, concentrically positioned in a carrier tube 14. Fixed within the charge tube 12 are shaped charges 16. Typically, the charge tube 12 is oriented in the carrier tube 14 such that the shaped charges 16 on each charge strip (not shown) align with weakened portions or scallops 18 formed in the carrier tube 14. A detonator cord 20 runs through suitable bores (not shown) in the perforating gun 10. Referring now to FIG. 1B, there is shown a cross-sectional view of the FIG. 1A perforating gun 10. Conventionally, the detonator cord 20 is coupled to the shaped charge 16 with an external clip 28. As is known, the shaped charge 16 may not detonate if the detonator cord 20 is not held firmly within a groove 30 formed in a post 31 extending from the rear of the shaped charge 16. Thus, in a typical assembly sequence, the shaped charges 16 are fixed within the charge tube 12 and then the detonator cord 20 is run through the grooves 30 of the shaped charges 16. Thereafter, the external clips 28 are snapped over the grooves 30 to retain the detonator cord 20 within the grooves 30. Next, the charge tube 12 is slid into the carrier tube 14. It should be noted that the annular clearance space 32 separating the charge tube 12 and the carrier tube 14 is relatively small in the region of the groove 30 and external clip 28. It should be noted that conventionally, the external clip 28 lies radially outward of the post 31 formed on the shaped charge 16. Disadvantageously, during installation of the charge tube 12 in the carrier tube 14, the external clip 28 sometimes hangs up with, gets caught by, and/or scrapes against the carrier tube 14 and becomes dislodged or damaged. If the external clip 28 is so damaged, the detonator cord 20 can fall out of the groove 30, in which case the shaped charge 16 may not detonate. An incomplete detonation of the shaped charges 16 can impair the effectiveness of the perforating activity and necessitate remedial work.\nThus, what is needed is a perforating gun that can be assembled without impairing the ballistic coupling of the detonator cord to the shaped charges. The present invention addresses these and other disadvantages of the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates in general to the field of hologram production and, more particularly, to applications of and interfaces for use with holograms.\n2. Description of the Related Art\nNumerous devices and techniques have been disclosed for producing and displaying various types of holograms including one-step techniques, updateable techniques, techniques relying on fringe pattern calculation, and traditional object/reference beam interference techniques. Moreover, the subject matter of the scenes displayed in such holograms varies from simple demonstration scenes, to artistic renderings, to portraits, to scenes of complex devices or spaces, and beyond.\nWith the development of various digital holography techniques, hologram subject matter can include virtually any computer graphics rendered scene, and can be based on numerous different data sources such as CAD information, physical measurements, photographic data, and the like. Holograms are particularly useful for visualizing complex three-dimensional spaces that users cannot normally interact with directly, such as large or small scale environments, dangerous environment, and inaccessible environments. The term “environment” is generally meant to include any sort of environment in which a person might find himself and/or representations of such environments, e.g., terrain, bodies of water, maps, buildings, campuses, facilities, vehicles, and other larger-scale physical spaces. Environments can also refer to smaller scales, the subject matter being visualized can include the inside of a device or apparatus, interior features biological specimens, or even sub-micron environments (e.g., the interior of a cell or molecule).\nDespite the fact that such environments can be displayed in holograms, the full potential of such visualization devices and techniques has not been realized for a number of reasons. Chief among those reasons are the absence of additional visualization tools and image enhancements to both ease the ability of a user to utilize what may be an unfamiliar visualization technique, and combine the environment visualization with related information.\nAccordingly, it is desirable to have improved systems and methods for producing, displaying, and interacting with holograms to overcome the above-identified deficiencies in the prior art."}
{"text": "1. Field of the Invention\nThe present invention relates generally to suturing devices and, more particularly, to an improved suturing gun which can be operated with one hand so as to facilitate the suturing operation.\n2. Description of the Prior Art\nIn the surgery field, mechanical suturing devices for suturing the skin are already known, and are generally produced in the shape of a gun. In use, a suturing gun of this kind is grasped by the surgeon with one hand, while he operates a forceps with the other hand in order to clamp the margins of the wounded skin together and maintain them in this clamped position so that the stitches delivered by the suturing gun can grip both the wound margins, such that the wound is properly closed. This procedure is very toilsome for the surgeon and often, to facilitate this operation, the mechanical suturing device is grasped by the right hand of the surgeon and the clamping of the skin margins is carried out by an assistant who operates a pair of toothed forceps to hold the wound margins clamped until the stitches have been applied. The latter solution facilitates the surgeon's task but requires the presence of two operators and therefore increases overall cost of the operation.\nIt is also known that during the suture operation the surgeon needs to have gauze tampons at his immediate disposal to plug the wound, as well as disinfectant liquids and suture adhesives in order to maintain the wound disinfected during the suture operation and sutured after the suture operation. The application of these products requires again the presence of a second operator during the suturing operation made by the surgeon, which operator must give the surgeon gauze tampons, disinfecting liquid and/or adhesive on request."}
{"text": "1. Field of the Invention\nThe present invention relates to peripheral device identification and interrupt management. More specifically, the present invention pertains to a method and system that enables device identification to be performed without a rigid time limit within which an initial interrupt from the device needs to be serviced.\n2. Related Art\nComputer systems and other electronic devices have become integral tools that are used extensively to perform a wide variety of useful operations in modern society. Applications of computer systems can be found in virtually all fields and disciplines, including but not limited to business, industry, scientific research, education and entertainment. For instance, computer systems are used to analyze financial data, to control industrial machinery, to model chemical molecules, to deliver classroom presentations and to generate special effects for movies. Moreover, computer systems along with other electronic devices are finding new applications as a result of advances in hardware technology and rapid development in software technology. The growing affordability of computer systems and electronic devices together with the abundance of useful new applications have fueled strong demand for such systems and devices.\nIncluded within this broad category of computers and electronic devices is the personal digital assistant (commonly referred to as a PDA). Specifically, as the components required to build a computer system have been greatly reduced in size, new categories of computer systems have emerged. One of these new categories of computer systems is the PDA. A PDA is a portable computer system which is small enough to be held conveniently and comfortably in the hand of its user. In particular, a popular size for the typical PDA approximates the size of a palm.\nThe PDA is usually a battery-powered device that is typically used as an electronic organizer having the capability to store and display a wide range of information which can include numerous addresses and telephone numbers of business and personal acquaintances, financial information, daily appointments, along with various other personal information. As such, the PDA is able to consolidate a wide variety of information and make the information easily accessible by its user. Therefore, PDAs are very useful and have gained wide popularity.\nSince PDAs are such small devices, full-sized keyboards are generally not efficient input devices. For instance, PDAs using keyboards usually have keyboard devices that are so small that a user typically cannot touch-type on them. Many PDAs thus employ a stylus and a digitizer pad as an input system. The stylus and digitizer pad work well for PDAs because the arrangement allows a user to hold the PDA in one hand while writing with the stylus onto the digitizer pad with the other hand. A small on-screen keyboard image can also be used for data entry. By eliminating the need for a keyboard, PDAs are very portable and can be carried along by their users wherever they go, even when on extended travel. On the other hand, due to their small size, PDAs usually have a modest set of built-in functions and it is often beneficial to couple peripheral devices therewith to expand functionality. The greater the variety of peripheral devices a PDA is capable of operating with, the broader the functionality it can achieve.\nMoreover, the management of communications port contention and power consumption affect the ease of use and the battery life of PDAs and other hand-held computers. A typical PDA includes limited communications port resources. For example, the Palm V personal organizer, manufactured by Palm Computing, Inc. of Santa Clara, Calif., has a single communications port available to interface with peripherals. When the communications port is exclusively occupied by a peripheral, the communications port is unavailable for other uses. Thus, a PDA with an active keyboard occupying the communications port, for instance, may be unable to perform an infrared (IR) synchronization process with another device.\nAnother ease of use issue is the degree of user intervention required to manage a communications port and identify peripheral devices coupled thereto. Requiring a user to manually open the communications port and identify a peripheral device, such as through the graphical user interface (GUI) or the buttons of the PDA or hand-held computer, is generally less preferable to automatically opening the communications port and identifying the peripheral device when the PDA receives information from a peripheral device upon coupling. Further, requiring a user to explicitly close the communications port is also generally less desirable than automatically closing the communications port after the peripheral device and the PDA no longer exchange data. In one existing solution, the PDA will automatically close a communications channel, which was opened to work with a keyboard peripheral, when the PDA is prompted with a signal initiating a synchronization process, such as when the PDA is placed in a cradle for a HotSync process. However, if a user forgets to explicitly close the communications port after using a keyboard that occupied the communications port, the user may be unable to perform an infrared synchronization process. In such cases, the synchronization process cannot be initiated until the situation is resolved. A user unfamiliar with the problem may even have the mistaken impression that the PDA has malfunctioned.\nFurthermore, requiring a user to manually close the communications port also affects power consumption. Battery powered systems, such as PDAs and other hand-held devices, are sensitive to applications that drain excessive power. Thus, PDAs mostly leave their communications ports in a low power standby mode, because an open communications channel may consume significant power. For example, the Palm V personal organizer has a communications port included in the processor that is shared between an infrared communications port and an RS-232 serial communications port. The Palm V hand-held organizer keeps these ports in standby mode to reduce power consumption. A user who forgets to close the communications port after a peripheral device is no longer being used may suffer a significantly decreased battery life, and be forced to recharge or change batteries often.\nEven if a user remembers to close the communications port when the attached peripheral device is no longer being used, the port has to be reopened for further communications with a peripheral device when use is resumed. When the communications port is reopened, the attached peripheral device needs to be identified. Identifying the peripheral device permits the PDA to open appropriate applications and select appropriate communication protocols. Some existing hand-held computers provide limited capability for identifying peripheral devices. For example, in an existing Palm organizer, this is implemented with two pins on a serial port. According to this technique, an interrupt is generated by bringing up one of the pins (e.g., making the pin “high”). An interrupt routine then checks the second pin to determine if it is high or low (e.g., high indicates modem, low indicates cradle). One problem with this technique is that it only allows two different types of devices to wake up the serial port. In other words, this peripheral device identification scheme is limited to identifying one of two devices, for example, a cradle versus a modem. As hand-held computer technology advances and its applications multiply, an increasing number of peripheral devices are becoming available for use with hand-held computers and PDAs, but the existing device identification technique are not capable of distinguishing among the numerous devices.\nSome other approaches to peripheral device identification are capable of identifying more than two types of devices. Typically, such approaches involve performing certain detection activities in a given time period immediately following the initial interrupt generated by the device as described above. According to these approaches, additional input is generated by the peripheral device and received by the hand-held computer during this prescribed time period. For example, the additional input may include one or more specific characters or additional interrupt(s). Based on the additional input, such as by monitoring for the specific inbound character(s) or by counting the additional interrupts, the hand-held computer determines the type of the peripheral device being coupled thereto.\nNevertheless, these identification schemes also have their shortcomings. In particular, these schemes impose a hard limit on the latency within which the initial interrupt generated by the peripheral device must be handled in order for the respective schemes to succeed. Such a time-critical interrupt response requirement is undesirable because the hand-held computer may be engaged in other computations when the initial interrupt is generated and thus may be unable to service the interrupt within the requisite time frame. Under these circumstances, peripheral device identification according to these time-critical schemes would fail.\nThus, in view of the foregoing problems that are inherent in existing peripheral device identification schemes, it would be highly advantageous to provide a peripheral device identification system and method that facilitates extendible identification of various peripheral devices, where the successful operation of the system and method does not require that the initial interrupt generated by the peripheral device be handled within a critical time period.\nAdditionally, a method and system for peripheral device identification should not require complete revamping of existing hand-held computer systems or PDAs. In other words, in implementing a viable method and system for peripheral device identification, components that are well known in the art and are compatible with existing hand-held computer systems need to be used so that the cost of realizing the method and system for peripheral device identification is low. By so doing, the need to incur costly expenditures for retrofitting existing hand-held computer systems or for building custom components is eliminated."}
{"text": "This invention relates to an amorphous or ultra fine-grained, diamond-like material that is substantially free of graphite and hydrogen, and is deposited in a film on a substrate in the form of nanometer-sized, tightly packed nodules of sp.sup.3 -bonded carbon, hereinafter referred to as \"nanophase diamond.\" In one method of preparing the material of the present invention, a laser beam is directed onto the substrate before and/or during deposition of diamond-like particles on the substrate to substantially eliminate graphite particles thereon. Laser energy focused on the substrate is helpful for preventing substantial build-up of graphite. Furthermore, laser energy on the substrate is also useful for preconditioning the substrate to facilitate bonding of diamond-like particles to the substrate.\nIn recent years, there has been great interest in producing a diamond-like carbon coating for a variety of reasons. First, diamond-like carbon has an extremely hard surface nearly impervious to physical abuse (abrasive or chemical) and is therefore quite useful as a protective surface. Diamond-like carbon is optically transparent (in, e.g., the infrared spectrum), and is therefore believed to be useful in a variety of optics applications such as protecting sensor optical circuits, quantum wells, etc. In addition, diamond-like carbon has been found to have a high electrical resistivity as well as high thermal conductivity--an unusual combination. Diamond-like carbon, when doped, can act as a semiconductor, thereby forming the basis of technology for microcircuitry that can operate under hostile conditions of high temperatures and high radiation levels. Therefore, great interest has been shown in developing techniques for obtaining diamond-like carbon films in commercial quantities for possible use in the semiconductor industry.\nWhile natural diamond is a generally well defined substance, so-called \"diamond-like carbon\" films are not well defined, possibly because of the many different methods of preparation which contribute unique aspects to the product. From a structural viewpoint, six allotropes of carbon have been identified--two for each of the numbers of dimensions through which the carbon atoms may bond. The two most important carbon allotropes of interest in the optics and semiconductor industries are the two dimensional sp.sup.2 bonding characteristic of graphite, and the three dimensional sp.sup.3 tetrahedral bonding which gives natural diamond its unique properties. The article Low Pressure, Metastable Growth of Diamond and Diamondlike Phases, by John C. Angus and Cliff C. Hayman, 41 Science pp. 1913-21 (August 1988), describes different types of diamond-like phases, methods of production, and possible uses, the disclosure of which is expressly incorporated herein by reference.\nNatural diamond has the highest hardness and elastic modulus of any natural material. Diamond is the least compressible natural substance known, has the highest thermal conductivity, and also has a low thermal expansion coefficient. Further, diamond is a wide band gap semiconductor, having a high breakdown voltage (10.sup.7 V/cm) in a saturation velocity (2.7.times.10.sup.7 cm/s) greater than silicon, gallium arsenide, or indium phosphide. Thus, diamond-like carbon films are potentially useful in the electronics industry as a protective, resistive coating with extremely desirable heat sink properties.\nSeveral different kinds of thin films having some of the properties of diamond can be made from pure carbon. Preferably prepared entirely without hydrogen, fluorine or any other catalysts, these materials can be distinguished by their internal structures. Three forms have been reported: (1) defected graphite, (2) i-C, and (3) \"amorphic\" diamond. See Collins et al., \"The Bonding of Protective Films of Amorphic Diamond to Titanium,\" 71 J. App. Phys. 3260 (1992), the disclosure of which is herein incorporated by reference. The term \"amorphic\" derives from a combination of the terms \"amorphous\" and \"ceramic,\" and has become a term of art to describe a material comprising amorphous carbon having diamond-like properties. In the first type of carbon-based material, diamond-like properties seem to accrue from a high density of defects in otherwise orderly graphite that trap electrons and shorten bonds to increase both strength and transparency. The other two materials, i-C and \"amorphic\" diamond, derive their enhanced characteristics from an abundance of carbon atoms linked by the sp.sup.3 bonds of diamond.\nCurrently, four major methods are being investigated for producing so-called \"diamond-like\" carbon films: (1) ion beam deposition; (2) chemical vapor deposition; (3) plasma enhanced chemical vapor depositions; and (4) sputter deposition.\nThe ion beam deposition method typically involves producing carbon ions by heating a filament and accelerating carbon ions to selected energies for deposit on a substrate in a high vacuum environment. Ion beam systems use differential pumping and mass separation techniques to reduce the level of impurities in the carbon ion fluence to the growing film. While films of diamond-like carbon having desirable properties can be obtained with such ion beam techniques, the films are expensive to produce and are achievable only at very slow rates of growth on the order of 50 angstroms per day to perhaps as high as a few hundred angstroms per day.\nThe chemical vapor deposition and plasma enhanced chemical vapor deposition methods are similar in operation and have associated problems. Both methods use the dissociation of organic vapors (such as CH.sub.3 OH, C.sub.2 H.sub.2, and CH.sub.3 OHCH.sub.3) to produce both carbon ions and neutral atoms of carbon for deposit on a substrate. Unfortunately, the collateral products of dissociation frequently contaminate the growing film. While both chemical vapor deposition and plasma enhanced chemical vapor deposition achieve film growth rates of practical levels, such films are of poor optical quality and unsuitable for most commercial uses. For example, if the films are amorphous, the internal stress is too high to permit growth to useful thickness; if crystalline, there is no bonding to the substrate. Further, epitaxial growth is simply not possible using chemical vapor deposition techniques.\nSputtering deposition usually includes two ion sources, one for sputtering carbon from a graphite source onto a substrate, and another ion source for breaking the unwanted graphite bonds in the growing film. In the typical sputtering method, an argon ion sputtering gun sputters pure carbon atoms off of a graphite target within a vacuum chamber, and the carbon atoms are condensed onto a substrate. Simultaneously, another argon ion source bombards the substrate to enhance the breakdown of the graphite bonding in favor of a diamond-like sp.sup.3 tetrahedral bond in the growing carbon film. The poor vacuum and relatively high pressure (10.sup.-5 to 10.sup.-4 torr) in sputtering deposition is cumbersome and tends to introduce contamination of the film on a level comparable to those encountered in chemical vapor deposition and plasma enhanced chemical vapor deposition.\nTherefore, while many attempts have been made to obtain high quality diamond-like carbon at commercial levels of production, the results have thus far been disappointing. The known methods recited above are deficient in many respects. While the ion beam deposition method produces a good quality film, its slow growth rates are impractical. The chemical vapor deposition and sputter methods are prone to contamination, yielding an unacceptable film in most circumstances. Additionally, the chemical vapor deposition technique creates a growth environment that is too hostile for fragile substrates. Indeed, the aggressive chemical environment characteristic of chemical vapor deposition could result in destruction of the very substrate that the diamond-like film is intended to protect. Moreover, all known methods require elevated temperatures, which often prove impractical if coating an optical substrate is desired. The known methods also involve complex and cumbersome devices to implement.\nChemical vapor deposition techniques have been useful in obtaining diamond-like carbon material denoted \"a-C:H,\" which is an amorphic carbon structure containing a significant amount of hydrogen. In fact, it is believed that hydrogen is necessary for permitting realistic growth of the metastable diamond-like material and that a reduction in the amount of hydrogen below about 20% degrades the film towards graphite. Of course the amount of hydrogen correlates with the proportions of sp.sup.3 to sp.sup.2 bonding and the correlation of diamond-like properties to graphite-like properties.\nIt appears that some ion-sputter techniques have been successful in producing a diamond-like carbon material without hydrogen--so-called \"a-C\" film. (See, Low Pressure, Metastable Growth of Diamond and Diamondlike Phases, page 920, supra). Such a-C diamond-like carbon is interesting in that the ratio of sp.sup.3 to sp.sup.2 bonding is believed to be high, with a corresponding increase in diamond-like properties. Unfortunately, the growth rates of such a-C materials are extremely slow and not commercially viable. In addition, internal stresses are too high to permit growth to useful thicknesses. Further, ion energies are too low to permit bonding with the substrates. Apparently, attempts to increase growth rates using increased power densities in such ion sputter techniques resulted in a decrease in the ratio of sp.sup.3 to sp.sup.2 bonding.\nProducing diamond-like carbon is just one example of the general problem of producing a layer of material having desirable physical properties where the material is extremely difficult to handle or manipulate. Examples of other such materials include semiconductors, such as silicon, germanium, gallium arsenide, and recently discovered superconducting materials which might be generally characterized as difficult to handle ceramics (e.g., yttrium-barium compounds). Therefore, it would be a significant advance to achieve a method and apparatus which could produce an optical quality diamond-like layer having an abundance of sp.sup.3 bonding (and therefore diamond-like qualities) in commercial quantities. Further, it would be significant if such method and apparatus were useful in producing layers of other types of materials which, using conventional technology, are difficult to handle or produce."}
{"text": "1. Field of the Invention\nThe present invention relates to sawmills and more particularly to a firewood cutting device.\nSince firewood is marketed by ricks it is desirous to provide a sawing device to cut logs in selected lengths which, after being split to firewood size, may be easily stacked to form the required dimensions for a rick of firewood.\nThis invention provides an apparatus for cutting a firewood forming log into a plurality of selected equal lengths to be thereafter split, as desired.\n2. Description of the Prior Art\nIt is well known to provide an equal length workpiece sawing machine comprising a pair of spaced circular saws to cut logs, planks or lath to equal length by passing the wooden material between the saws.\nU.S. Pat. No. 1,123,056 discloses a cord wood cutting device having a pair of frame journalled power driven saws including a disk journalled by the frame parallel with the saws and provided with radially projecting stops which receive firewood logs and move the logs toward and beyond the saws as the logs are successively positioned on the stops. The cord length sawed logs fall by gravity to a conveyor for movement toward a disposal point.\nThis invention is distinctive over this patent by providing a manually operated slip clutch flat belt drive for a conveyor feeding logs horizontally toward a plurality of equally spaced saws in a plane disposed below a horizontal plane taken diametrically through the saws wherein the slip belt control permits interruption of the feed conveyor action in the event of malfunction of any of the equipment or wood being cut."}
{"text": "The present invention relates to a semiconductor storage element, a semiconductor device and manufacturing methods therefor as well as a semiconductor device in which a semiconductor storage element and a semiconductor switching element are consolidated on an identical substrate and a manufacturing method therefor. The present invention also relates to portable electronic equipment and an IC card provided with the above-stated semiconductor device.\nConventionally, there has been a flash memory as a memory device being electrically writable and erasable. As shown in a sectional view of FIG. 19, the conventional flash memory has a floating gate 906 made of polysilicon and arranged on a semiconductor substrate 901 via a first oxide 904, a control gate 907 made of polysilicon and arranged on the floating gate 906 via a second oxide 905, and a pair of source/drain diffusion regions 902 and 903 arranged in the surface of the semiconductor substrate 901. This flash memory constitutes a field-effect transistor (hereinafter referred to as FET), where the control gate 907 operates similarly to the gate electrode of the FET. The first oxide 904, the floating gate 906 and the second oxide 905 are arranged between the control gate 907 and the semiconductor substrate 901. That is, the flash memory includes a memory film (i.e. the floating gate 906) having a charge storage function and arranged in a position corresponding to the gate insulator of the FET, so that the flash memory executes storage, read and erase operations by changing the threshold voltage of the FET according to the amount of charge accumulated in this memory film.\nThe conventional flash memory has a problem that the memory characteristics deteriorate due to so-called overerase (Masuoka Fujio, “Flash Memory Technology Handbook”, Science Forum Co., Ltd., Aug. 15, 1993, P 55-58). That is, the overerase is caused by excessive extraction of electrons or the injection of positive holes when the flash memory executes the erase operation by extracting the electrons accumulated in the floating gate 906 or injecting positive holes into the floating gate 906. Due to this overerase, the floating gate 906 is electrically positively charged, and the FET is turned on by the influence of the positive electric charge, disadvantageously causing a leakage current flowing between the source/drain diffusion regions. This leakage current is attributed to the fact that the control gate 907 corresponding to the gate electrode of the FET and the floating gate 906 that serves as the memory film are stacked and the FET is turned on and off only by the electric charge retained in the floating gate 906.\nA memory cell array formed of the aforementioned flash memory cells has a problem that the current of a selected memory cell cannot be extracted during the read of the selected memory cell due to the leakage current generated from non-selected memory cells, causing a read error."}
{"text": "Computerized social networks have a benefit that is derived from large size. Larger sized networks, such as Facebook, have a significant number of users making a significant number of interactions with other users and content items on the network to provide the operators of the network with social data that can reasonably approach a point of meaningful information that can then be utilized to present to its users recommendations, item listings, and other information based on how other users in the network have behaved. Smaller networks and e-commerce systems typically do not have the volume of users or user interactions to effectively generate their own meaningful data. Additionally, the topology, homophilous nature, and dimensionality of the structure of the users and network may be insufficient for a desired social signal. Further, real users can provide additional problems, such as difficulty in building a critical mass, scaling problems with too many users, users not interacting with consistent regularity, too many similar users (which results in something like over-fitting in a machine learning application), and/or other problems."}
{"text": "This application is related to U.S. application Ser. No. 924,775 and 924,639.\nThe invention concerns belts or straps, more particularly elevator conveyor belts. It relates to a method and a device for splicing the ends of belts, either to increase their length, or to make them endless by joining the belt to itself in a manner which is reliable during operation. The fastening device as well as the process will be described below, as an example, for application to an elevator belt, with the understanding that they are not limited to this use.\nOne conventional method of splicing the ends of belts is hot splicing after the ends of the reinforcing material have been interlaced.\nVarious forms of these hot splicing techniques are described in a number of patents, such as DE No. 1,165,354 of Franz Clouth Rheinische Gummiwarenfabrik; FR No. 74.03141 and DE No. 907,996 of Continental Gummi-Werke; FR No. 1,395,634, FR No. 1,582,190 and FR No. 1,440,605 of Pneumatiques, Caoutchouc Manufacture et Plastiques Kleber Colombes; and U.S. Pat. No. 173,686 of Goodyear Tire and Rubber. One method is recommended in German standard DIN 22131. Another method is described by Mr. Gozdiff of Goodyear in a paper entitled \"Factors relating to vulcanized splice reliability for steel cable reinforced conveyor belting\", delivered to the 125th Meeting of the Rubber Division, American Chemical Society, in Indianapolis on May 8-11, 1984. Finally, an article by H. P. Lachmann entitled \"A survey of present-day conveyor belt technology\", published in Bulk Solids Handling volume 4, number 4, December 1984, reviews the different technologies that can be used. Other proposals have been described in, for example, U.S. Pat. Nos. 2,446,311 and 3,093,005.\nExamples of the prior art in hot splicing are illustrated in FIG. 1. There are, however, conditions which make hot splicing techniques inapplicable: for example the length of such a splice may be incompatible with the space available in the sheath or alongside the elevator, or again the mechanical strength may become insufficient to guarantee trouble-free operation of the elevator if the temperature of the products transported or that of the gases circulating in the sheath is greater than 100.degree. C.\nSpecifically, in the case of interlacing of metal cords, the two ends of the belt(s) are bonded together by the rubber mixture separating the ends of the cable(s). It is known, however, that as the temperature increases, the mechanical properties of elastomer-based mixtures decrease; the same applies to the bonding forces between the rubber and the metal. As a consequence, the tensile strength of such a splice decreases as the temperature rises.\nTaking into account the risks encountered with a splice using hot adhesion and interlacing of metal cords, handling engineers have suggested replacing the adhesive bond with a mechanical bond designed to clamp the two ends of the belts(s) against one another. These techniques are referred to as \"fastening\", and are described, for example, in French Pat. No. 1,320,222 of Pneumatiques, Caoutchouc Manufacture et Plastiques Kleber Colombes, or in advertising materials of specialized companies such as Goro or Flexco.\nOther mechanical fasteners are described in, for example, U.S. Pat. Nos. 2,447,855 and 1,918,255.\nThe compression force is exerted by metal flanges which are passed through by clamping bolts. The principle applied in mechanically splicing the ends of a belt or belts is theoretically more satisfactory than the hot-adhesion process, but an analysis of phenomena associated with operation of this type of device shows that there is only a slight improvement in operating reliability. This is because the lateral plates have a tendency to move away from one another under the tensile force exerted on the two ends of the belt. To remedy this problem, two lines of bolts are generally used to clamp the plates, with the line of bolts placed closest to the tension zone being designed to limit movement of the plates.\nFIG. 2 illustrates conventional fastening techniques such as described in U.S. Pat. Nos. 4,489,827, 1,803,354, 2,330,434, 2,330,435, 3,618,384, and 4,450,389 and provide an example of changes which occur during operation, in terms of distribution of pressure over the belt ends.\nWhen the fastening device is initially clamped, pressure distributed uniformly over the entire extent of the two clamped ends. When operation begins, the lateral parts are displaced, which tends to decrease the pressure in zone X and increase that in zone Y. As a result of the increased pressure, the rubber mixture located in zone Y tends to be expelled and to creep, i.e. rubber is displaced from the most highly compressed zones towards those least compressed. The effect of this creep is to encourage the clamping plates to move closer together in zone Y, which accentuates the effect.\nAs confirmation of this analysis, it is commonly observed that the second line of bolts has completely loosened, which proves that the pressure effect exerted on the ends of the belt(s) is due not to the bolts but to a rotary movement of the clamping plates. Since the clamping force in zone X has decreased considerably, the strength of such a clip fastener consists only of the retention of the metal cords in zone Y and the frictional forces existing in zone X between the belt and each of the clamping plates. When zone Y is unclamped, an examination of the ends of the belt(s) in said zone shows that the tensile stress has been so high that there has been local destruction of the rubber mixture and of the bond between the rubber and the metal cords. This fact makes the strength of such a clip fastener very problematical, especially when this mechanical effect is combined with the effect of temperature, since it is well known that increased temperature accelerates and facilitates the creep of elastomer-based mixtures, and decreases the strength of the bond between said mixtures and steel cords.\nTo remedy this problem, it is possible to attempt to increase the pressure exerted by the lateral plates by locking the ends of the metal cords.\nPatent DE No. 2,341,992 of Bernhard Beumer Maschinenfabrik describes such a solution, in which each metal cord is stripped at its end of its rubber covering, and said end is inserted into a clamping device using screws. Such a technique, time-consuming and difficult to implement, presents a further risk due to the design of the metal cord clamping zone: if the clamping pressure is not properly controlled, there is a definite risk of cutting the metal cord, which would nullify the anticipated effect. In addition, the screws have a tendency to loosen under the action of vibration and temperature, and therefore require constant monitoring.\nA similar solution of the type described above using bores steel balls to lock the ends of metallic cables is also proposed in U.S. Pat. No. 3,105,390.\nA different solution, used in particular to lock pretensioning cords in prestressed structures, involves stripping the ends of the metal cords, unstranding them, i.e. untwisting the constituent strands to spread out the end, and pouring around it a metal with a low melting point. This technique is highly reliable when it can be used, but pouring the molten metal requires that the clip fastener be placed in a vertical position, which implies either that one of the drums of the elevator can be moved - which is not always possible--or that the clip fastener can be placed at the top of the elevator, although pouring molten metal onto the ends of metal cords at a height of several dozen meters is tricky and even dangerous. In addition, such a device is practically non-removable, which does not facilitate maintenance of the elevator or replacement of the belt.\nThe object of the invention is therefore a fastening device for splicing belts, and more particularly belts for bucket elevators, which comprises a reinforcement made of manmade fabric or metal cords. Other objects of the invention are the process for creating splices using the device described, and application to an elevator conveyor belt.\nThe fastening device which is the object of the invention makes it possible to create, using a method that is as simple as conventional splices and is completely safe since there is no handling of hot products or adhesives or heating equipment, a splice between two adjacent ends of two belts to produce a longer belt, or between the two ends of a single belt to produce an \"endless\" item, ready for operation on a conveyor or elevator.\nThe device, which requires only simple preparation of the ends, acts on the one hand by locking by means of a loop created around a locking core, and on the other hand by means of mechanical clamping of said ends over the greatest possible length.\nThis fastening device comprises, on the one hand, two central plates, designed solely for installation and serving to support the ends of the belts(s), to which are attached, on the upper part, two installation wedges allowing the assembly to be aligned, and on the other hand, two locking cores making it possible to create loops in the ends of the belt(s) and, finally, two lateral clamping parts each fitted with a ball joint.\nThis device acts like a pivoting self-clamping clip by means of its various component parts, with the lateral clamping parts acting to clamp the belt by means of through bolts.\nThe characteristics and variants of the invention will be better understood by reading the description below taken in conjunction with the accompanying."}
{"text": "1. Field of the Invention\nThe present invention relates to a write-once recording medium comprising a management information area for recording management information for managing a recorded state and a user data area for recording user data, and a method and apparatus for recording data onto the write-once recording medium, and a method and apparatus for reproducing data from the write-once recording medium.\n2. Description of the Related Art\nOptical discs are information recording media that are compatible between a plurality of apparatuses. Examples of optical discs include, for example, the CD which was originally developed for recording music and recently are widely used for recording data, DVD which was developed for recording digital video and are rapidly becoming widespread, and next-generation optical discs which are currently being developed for recording high-quality video, such as for high-definition television or the like.\nThese optical discs are roughly categorized into three groups in terms of the recording/reproduction method: a read-only type; a write-once type; and a rewritable type. Write-once optical discs are made of a recording material, from which recorded data cannot be erased. Data can be recorded once in the same place on the medium. For example, CD-R and DVD-R are write-once optical discs. Write-once optical discs are also referred to as write-once-many-read optical discs. For rewritable optical discs, data can be recorded many times in the same place on the medium. For example, CD-RW and DVD-RAM are rewritable optical discs.\nWrite-once optical discs are not rewritable due to the characteristics of the recording material. Therefore, it is necessary to manage recorded areas and unrecorded areas. In a CD-R, the recording start position and the last recorded position of each music composition are managed so as to incrementally record in units of a music composition. The unit of a music composition is called a track (hereinafter also referred to as a “TRACK” so as to distinguish it from a track indicating a groove). In a CD-R, data is recorded sequentially in units of a track from its inner periphery to its outer periphery. A plurality of recorded TRACKs are managed as a session including a plurality of recorded TRACKs (hereinafter referred to as a “SESSION”). An area for storing management information for managing a recorded state is called a program management area (hereinafter also referred to as a “PMA”).\nFor DVD-R, an RZone (corresponding to a TRACK of CD), a Border (corresponding to a SESSION of CD), and a recording management area (hereinafter also referred to as an “RMA”; corresponding to a PMA of CD) are defined (for example, SFF 8090i standard, “Mt. Fuji Command for Multimedia Devices”, Section 4.16 entitled “Recording for DVD-R media”).\nThere is a technique for applying a method for managing areas allocated on DVD-R compatibly to read-only DVD-ROM (for example, Japanese Laid-Open Publication No. 2001-148166).\nIn a data structure recorded in an RMA of DVD-R, a recording management data (hereinafter also referred to as an “RMD”) is defined.\nFIG. 23 shows a field 3 of an RMD. The field 3 of an RMD stores information indicating the positions of a maximum of 512 Borders (corresponding to a SESSION of CD).\nFIG. 24 shows a field 4 of an RMD. FIG. 25 shows fields 5 to 12 of an RMD. The field 4, and the fields 5 to 12 of an RMD store information indicating the positions of a maximum of 2302 RZone's (corresponding to a TRACK of CD).\nHowever, in conventional techniques for managing recording areas, the number of recording areas to be managed is currently already large. If the conventional techniques are applied without modification to a larger capacity of future media, management information is adversely enlarged."}
{"text": "The present relates generally to bearings. More particularly, the present invention concerns segmental bearings of the type used with a rotating machine of substantial size.\nStringent requirements are placed on shaft bearings of power plant turbosets and power plant generators as a result of the great rotating weight associated with rotors, shafts and the components supported thereby. By virtue of the high static and dynamic force loads and the external effects generated by the working medium, these shafts bearings are subjected to high specific bearing surface compressions (20-60 bar) and journal circumferential velocities which substantially exceed 100 m/sec.\nA hydrodynamic bearing is used exclusively for supporting such shafts has proven to be thoroughly satisfactory for the operating conditions. In addition, the hydrodynamic bearing has attained a high state of development. With a hydrodynamic bearing it is possible to produce a sufficiently thick lubricating film by self-lubrication. The required thickness for such a lubricating film is determined by the magnitude of the deflection under load that is experienced by the guide block in the case of an axial bearing such as used in a hydroelectric generator or by deflection of a bearing box segment in a radial bearing of a power plant turboset. In these hydrodynamic bearings, deflection of the bearing surface of the guide block on bearing box segment is caused by hydrostatic pressure in the oil lubricant film created during rotation of the shaft. Such an hydrostatic pressure acts as a stress nonuniformly distributed over the entire bearing surface with the stress decreasing toward the edges of the respective guide block or bearing box segment and causes a bending to occur in the bearing surface.\nThis mechanical deflection caused by hydrostatic pressure has superimposed thereon deflections which result from the heating of the guide block or bearing box segments and the associated supporting elements. It is to be understood that supporting plates are exposed to this thermally-induced deflection in the case of axial bearings whereas supporting segments are exposed to the thermally-induced deflection in the case of radial bearings.\nIn order to avoid removal of the lubricant film and subsequent metallic contact of the bearing surfaces, sometimes referred to as mixed friction, the thickness of the lubricating film which occurs during operation must always be greater than the resulting deflection produced by the combined stresses: i.e., induced by the hydrostatic pressure and the thermally induced deflection. It is only when such a thickness condition is satisfied that the bearing surface and the cooperating shaft surface will remain separated from one another by the lubricating film. Without such separation, damage can occur which can eventially lead to destruction of bearing surface and adjacent shaft.\nOne object of the present invention consists of providing a bearing of the type suitable for use with large rotating machines in which the bearing surfaces deflect in such a manner that the bearing surface remains essentially planar and moves essentially parallel to the original undeflected position. These deflections are caused by the mechanical and thermal stresses noted above. Moreover, it is an object of the present invention to confine the geometrical deviations of the deflected bearing surfaces so that those deviations lie within very narrow limits from their unstressed configuration and so that the deviations are less than the lubricating film thicknesses conventionally occurring during operation. Accordingly, liquid friction is guaranteed under all operating conditions."}
{"text": "Conventional carding apparatus used in the manufacture of textile yarn from short staple fibres comprises a cylinder having a surface which carries an array of carding points. As the cylinder rotates, the carding surface picks up fibres and carries them into engagement with a series of carding flats which are also provided with carding points. The carding flats are moved in the same direction as that part of the carding surface of the cylinder with which they engage, but at a lower speed. Fibres carried by the cylinder are therefore combed by the carding surfaces of the cylinder and flats, so that impurities or tangles in the fibre are removed on to the flats. The flats themselves are driven along an endless path and, on disengaging from the carding cylinder are driven past a brush which frees the surfaces of the flats from the waste material.\nHitherto, the surfaces of carding flats have been composed of a flexible web, which is secured to a rigid substrate, the points being in the form of staples projecting through the web from the side fixed to the substrate. Alternatively, the carding surfaces of the flats have been provided by metal strips having saw-tooth edges secured in the surface of a substrate. Both these constructions have suffered from the disadvantage that the fibres removed from the surface of the cylinder are difficult to separate from the carding flat because in the former case the points are of uniform cross section and in the latter case the saw tooth edges inevitably include burrs which trap the fibres."}
{"text": "The disclosure generally relates to small containers or boxes, such as those used for displaying or packaging jewelry and other items. Typically, jewelry, watches, and other such items are displayed on a store rack, which may be placed on a counter. Exemplary racks used for display may include one or more flat display faces that may be mounted on a rotatable base to enable a customer to quickly look through the items on display.\nEach display face of a typical display rack may include a plurality of regularly spaced hooks or openings, each of which is configured to support a particular item for display. Items are usually placed within boxes that are either open or transparent to allow the customers to view the contents. Open boxes in particular are used to display jewelry and other items that have visual appeal. To support the boxes on the rack, disposable attachments to the boxes, such as removable plastic hooks, have been used in the past. For example, a removable plastic hook that attaches into an opening of the box may be used to hang the box from an opening in the display rack. When an item in its box is selected by a customer, a sales associate may remove the box from the display, remove and dispose of the hook, and retrieve a lid or other closure of the box such that the item may be purchased in its packaging, which may be gift wrapped.\nThis type of known product display and packaging arrangement presents certain disadvantages. For one, the cost of merchandise packaging is increased by the use of the disposable plastic clips. Further, waste is created by the disposable clips used to mount the open boxes to the display. Even further, the lids for the open display boxes are stored near the display for use when items are selected for purchase, thus taking up space and time for the sales associate to retrieve them."}
{"text": "The rotary connector is used as electrical connection means of, for example, an air bag system. The rotary connector includes a rotator fixed to a steering wheel mounted with an air bag, and a stator rotatably supporting the rotator, and accommodates a flexible flat cable in a wound state in an internal space formed by the rotator and the stator. With this configuration, even when the rotator rotates with the rotation of the steering wheel, the electrical connection state can be maintained.\nIn a fixing structure for fixing the rotary connector to the combination switch, locking is used so that fixing work can be performed more easily as compared to the case of using a screw. That is, as shown in FIG. 17, projecting pieces 103 are formed on a stator 102 of a rotary connector 101, and as shown in FIG. 18, a locked portion 104 to which the projecting piece 103 is locked is formed in a bracket of a combination switch.\nThe projecting piece 103 has a rod shape rectangular in cross section, and at the tip portion of the projecting piece 103, a locking claw 105 is formed which protrudes laterally, relatively moves with respect to the locked portion 104, and is locked to the locked portion 104 when passing over the locked portion 104.\nAt the time of the fixing, the projecting piece 103 is elastically deformed in a stage before the locking claw 105 is locked to the locked portion 104, and after the locking claw 105 has passed over the locked portion 104, the projecting piece 103 elastically returns to the normal state, whereby the locking claw 105 is locked to the opposing side of the locked portion 104.\nHowever, in such a locking structure, in order to return the elastically deformed projecting piece 103 to the original shape, a certain clearance needs to be provided between the locking claw 105 and the locked portion 104 facing the locking claw 105. Hence there is a drawback that looseness occurs in a direction in which the projecting piece 103 is inserted.\nWhen such looseness occurs, an abnormal sound tends to occur. Moreover, as shown in FIG. 19, the rotary connector 101 is provided between a steering wheel 106 and a combination switch 107, and since a harness for electrical connection is provided in the rotary connector 101 on each of the steering wheel side and the steering column side, the harness may push up the rotary connector 101 depending on a wiring state of the harness (not shown) as indicated by an arrow in FIG. 19. In this case, the pushed-up rotary connector 101 may contact the steering wheel 106 to cause generation of the abnormal sound and other inconveniences.\nFor preventing the looseness of the rotary connector, a technique such as Patent Document 1 described below has been disclosed. A mounting structure of this technique includes a displacement absorbing member for preventing looseness, and an elastic piece or a cushioning material elastically striking the combination switch is provided on a stator.\nWith such a configuration, the displacement absorbing member brings the rotary connector into a pushed-up state, so that the rotary connector can be prevented from being lifted unexpectedly. However, a space is required between the rotary connector and the combination switch, and furthermore, extra materials and members are required. This makes it impossible to reduce the size and cost. Especially, when the cushioning material is provided, deterioration in the material and dropout of the cushioning material are considered, and it is difficult to maintain the effect."}
{"text": "For the coating of a moving base, such as a paper board web, with a coating agent, various size presses are commonly used. In size presses, the paper or board web is passed through a nip formed by size press rolls, in which nip the coating agent is transferred onto the face of the paper or board web.\nIn the prior art, a number of size presses of different types are known, of which one application, which has been in use for a long time, is one in which the rolls in the size press are arranged substantially on the same horizontal level so that the paper or board web runs through the size press nip substantially vertically from the top downwards. In such an application, as a rule, in the nip between the size press rolls, there is a size pond, through which the paper or board web runs. When the web passes through the pond, the size adheres to both sizes of the web, whereupon, in the nip, the size press rolls spread and smooth layers of suitable thickness onto the web faces.\nA second embodiment of a size press is one in which, by means of some suitable coating means, such as bar or blade coaters, size films of suitable thickness are first spread onto the faces of the size press rolls, from which said size films are then transferred onto the web faces in the roll nip. A third prior art embodiment of a size press is a so-called gate roll size press, which is, in two-sided surface sizing of a paper, composed of size rolls, which are in nip contact with each other. A gate roll size press first comprises applicator rolls proper, the web being passed through the nip between said rolls. Size films are spread onto the faces of said applicator rolls by means of gate rolls, which comprise a metering roll and a transfer roll. The transfer roll is in nip contact both with the metering roll and with the applicator roll proper. The coating agent is fed into the nips between the gate rolls, i.e. between the metering rolls and the transfer rolls, where it forms a pond. From both ends of the ponds, as a rule, an overflow is provided to the return circuit. By the effect of hydrodynamic forces, coating agent passes through the nip between the metering rolls and the transfer rolls, forming size films both on the metering rolls and on the transfer rolls. The size film placed on the transfer roll is transferred into the nip between the transfer roll and the applicator roll, which nip smooth and thins the film onto the applicator roll. In the nip between the applicator rolls, the film are transferred further onto the web.\nThe present invention is expressly related to the size presses of the last-mentioned gate roll type. The construction of conventional gate roll size presses is usually such that one of the applicator rolls is journalled fixedly on the size press frame, whereas the other applicator roll is arranged displaceable so that its bearing bracket is linked pivotally to the frame of the size press. The transfer roll that is in nip contact with the fixed applicator roll is linked pivotally to the size press frame by the intermediate of a loading arm and, further, the metering roll that is in nip contact with said transfer roll is linked pivotally to the size press frame by the intermediate of a loading arm and, further, the metering roll that is in nip contact with said transfer roll is linked pivotally to the bearing bracket or loading arm of the transfer roll by the intermediate of its own loading arm. In a corresponding way, in conventional solutions, the transfer roll that forms a nip with the displaceable applicator roll is linked pivotally to the bearing bracket or loading arm of the applicator roll by the intermediate of its own loading arm or by its bearing bracket and, further, the metering roll that is in nip contact with said transfer roll is linked pivotally to the bearing bracket or loading arm of the transfer roll. On the other hand, the loading cylinders of each roll are linked to the loading arm of search roll and, at the opposite end, to the machine member to which the roll concerned is linked pivotally.\nIn wide and high-speed paper machines in particular, the prior art solutions are associated with considerable problems of vibration of the gate rolls, resulting from the mode of linkage of conventional gate rolls. For example, the metering roll at the side of the displaceable applicator roll, i.e. the outermost roll in the system, is fixed to the frame by the intermediate of a chain formed by three separate articulated joints. By means of the construction concerned, it is very difficult to provide a sufficiently rigid fastening of the rolls in view of vibrations. Besides wear of the constructions of the frame, the vibrations also cause quality defects in the coating band an increased wear of the roll coatings. The wear of the roll coatings is one of the most important problems of the gate roll size presses."}
{"text": "1. Technical Field\nThe present invention relates to a liquid ejecting apparatus such as an ink jet type recording apparatus and to a liquid ejecting apparatus in which a liquid stored in a liquid storage member is introduced into a pressure chamber and the liquid introduced into the pressure chamber is ejected from a nozzle as liquid droplets.\n2. Related Art\nA liquid ejecting apparatus is an apparatus that includes a liquid ejecting head and ejects various types of liquids from a liquid ejecting head. As the liquid ejecting apparatus, for example, there is an image recording apparatus such as an ink jet type printer (hereinafter, simply referred to as a printer) or an ink jet type plotter.\nThe liquid ejecting head has been developed which uses a cartridge type liquid storage member, because of ease of handling and distribution. For example, the printer has been widely prevalent which uses an ink cartridge in which a liquid-phase ink is sealed.\nFIG. 5 is a schematic view illustrating an ink cartridge 3 as a liquid storage member in which a liquid-phase ink is sealed and a recording head 2 as a type of a liquid ejecting head.\nWhen the ink cartridge 3 is provided in the recording head 2, the ink inside the ink cartridge 3 is introduced into the recording head 2 side via an ink introduction hole 32 opened to a front end side of an ink introduction needle 92 by inserting the ink introduction needle 92 of the recording head 2 into the ink cartridge 3.\nThe ink introduced into the recording head 2 is introduced into a common liquid chamber 96 (also referred to as a reservoir or a manifold) via an introduction path 95 inside the recording head 2. The ink introduced in the common liquid chamber 96 is supplied to each of a plurality of pressure chambers 98 communicating with the common liquid chamber 96. Then, the pressure inside the pressure chamber 98 is changed by driving a piezoelectric vibrator, a heating element or the like that is a type of a pressure generation unit and the ink is ejected from a nozzle 99 via the pressure chamber 98 by controlling the change in the pressure.\nHere, when exchanging the ink cartridge 3 (when the ink cartridge is inserted into the ink introduction needle), there is a case where an excessive positive pressure or an excessive negative pressure may occur inside an ink flow path of the recording head 2 by button operation of an equipment user, change in a relative position occurred between the ink cartridge and the ink introduction needle, impact from the outside, or the like. Then, when the pressure occurred at this time is transmitted to the nozzle 99, a meniscus formed on the nozzle 99 is destroyed.\nFIGS. 6A to 6C are enlarged views of the vicinity of the nozzle 99 illustrated in the dashed circle in FIG. 5 and schematic views describing states of a meniscus M in the nozzle 99 when the change in the pressure is transmitted.\nA state illustrated in FIG. 6A is an appropriate state, however, there are cases that the meniscus M is excessively drawn into the pressure chamber 98 side from the inner peripheral surface of the nozzle 99 as illustrated in FIG. 6B and the meniscus M is inflated outwards from an opening surface of the ejection side of the nozzle 99 as illustrated in FIG. 6C. Then, when the pressure in the meniscus M exceeds a withstanding pressure, the meniscus M is not appropriately formed (a destroyed state), the ink may not be ejected, conversely, the ink may be leaked, and then dot failure may occur.\nIn order to solve such a problem, for example, a differential pressure valve structure has been known in the related art in which a valve is provided inside the ink flow path and the valve is open and closed according to the pressure inside the ink flow path, as disclosed in JP-A-2008-149646.\nThe structure disclosed in JP-A-2008-149646 is a valve that is configured such that an object of which is to stably supply the liquid to the downstream side by opening the valve in a case where the amount of the liquid at the downstream side is reduced by the liquid normally ejected from the nozzle and the downstream side has a negative pressure. However, there is a case where the valve does not respond to a sudden change in pressure from the impact outside. Thus, when the sudden change in the pressure occurs, the transmission of the change in the pressure to the nozzle may not be prevented and it is difficult to prevent the destruction of the meniscus."}
{"text": "The present invention relates to a spindle device in which the spindle bearings are lubricated with an oil jet.\nOil jet lubrication has proved superior to oil mist lubrication for very high speed bearings, cooling and lubricating them much better than oil mist which is now generally used for lubricating very high speed precision ball bearings in machine tools or others.\nExperimental oil jet lubrication has been successful in laboratories, but the practical use was not yet realized because of problems in expensive oil consumption and environmental oiling."}
{"text": "The world has recognized the critical need to decouple economic growth from resource impact. In particular, Europe is aimed at increasing industrial competitiveness whilst drastically reducing resource and energy inefficiencies. The underlying principle is to develop enabling technologies and solutions along the value chain to “do more with less”.\nThe following objectives have been proposed:                1. A reduction in fossil-fuel energy intensity of up to 30% from current levels by 2030 through a combination of, for example cogeneration-heat-power, process intensification, introduction of novel energy-saving processes, and progressive introduction of renewable energy sources within the process cycle.        2. By 2030, up to 20% reduction in non-renewable, primary raw material intensity versus current levels, by increasing chemical and physical transformation yields and/or using secondary and renewable raw materials with proven sustainability advantages.        \nThe possibility to burn glycerol obtained as a by-product in the synthesis of biodiesel is an area where these principles are most relevant since biodiesel along with bioethanol is currently the major biofuel in the market and, in addition, its manufacture is inefficient because not all the oil feedstock is converted into biofuel. The concomitant by-product glycerol is obtained along with fatty acid methyl esters which are the current automotive biofuel.\nMany attempts have been conducted in order to burn glycerol obtained as a by-product in the synthesis of biodiesel as a way to obtain renewable energy for burning purposes. However, crude glycerol display unfavourable fuel characteristics and it does not burn consistently in commercial burners.\nAgirre I et al. “Glycerol acetals, kinetic study of the reaction between glycerol and formaldehyde”, Biomass and Bioenergy, Pergamon, Amsterdam, NL, Vol. 35, no. 8, pp. 363-3642, discloses a composition resulting from the reaction between glycerol and formaldehyde. The study is aimed at determining kinetic parameters of this reaction by using ionic resins as catalysts.\nRuiz V R et al. “Gold catalysts and solid catalysts for biomass transformations: Valorization of glycerol and glycerol-water mixtures through formation of cyclic acetals”, Journal of Catalysys, Academic Press, Duluth, Minn., US, Vol. 271, no. 2, pp. 351-357, discloses a composition resulting from the reaction between glycerol and formaldehyde. In this case the study is aimed at discovering new catalytic systems (Gold catalysts).\nIn these both documents the low conversions described must result in compositions that contain, on stoichiometric basis, large quantities (>100 ppm) of formaldehyde.\nTherefore unreacted formaldehyde in the form of formaldehyde or a formaldehyde precursor, trioxane, is present in the compositions described therein. This presence of formaldehyde is absolutely not desired since it is well known in the art that, even at low concentrations, is a severe irritating substance, being even carcinogenic to humans after a long-term exposure.\nTherefore, although these documents disclose a composition including glycerol and glycerol formal which may be used as a biofuel ingredient, the current industrial standards would never use this toxic composition (including formaldehyde) because the authorities would not approve this product under current regulations\nAll these problems have prompted the inventors to develop a novel biofuel composition with improved burning characteristics when compared to crude glycerol which fulfill the need for efficient and renewable fuels and is free of toxic components, such as formaldehyde.\nIt is therefore an object of the present invention to provide a novel non-toxic glycerol-based biofuel formulation which circumvents the inability of crude glycerol to burn consistently in current burners. It is also an object of the invention to provide a procedure to prepare said biofuel. It is a further object of the present invention to show the better burning characteristics of the biofuel described in the present invention as compared to crude glycerol."}
{"text": "In recent years, digital cameras and personal computers have been widespread, and data generated by the digital cameras or the personal computers can be printed by inkjet printers or laser printers. These printers perform printing by adhering ink or toner on a recording medium, such as paper. The ink or the toner is stored in a cartridge or the like and supplied to the printers. The cartridge is replaced when the ink or the toner is consumed. In the case of large ink or toner consumption, as the frequency of replacement of the cartridge is increased, the cost for purchasing the cartridge becomes higher.\nThere have been proposed methods for calibrating data to be printed so as to reduce ink or toner consumption without deteriorating the quality of the data.\nFor example, according to a method disclosed in US 2005/0063749, data is calibrated according to an ink reduction amount set by a user and then transmitted to a printer driver, thereby reducing ink consumption in a printer.\nHowever, even with the same ink reduction amount, the degree with which quality of printed image is deteriorated as a result of calibration varies depending on the type of image to be printed based on data (photographs taken by digital cameras or documents, for example) or the type of recording paper (glossy paper or plain paper, for example).\nFor example, when a large ink reduction amount is set for data for photographs or the like, a resultant image quality may be deteriorated more than expected. This may necessitate a user to reprint the same data with less ink reduction amount, ending up with larger ink consumption. On the contrary, when a small ink reduction amount is set for data for a document or the like, ink consumption may not be reduced effectively. In this manner, it is difficult for a user to set an appropriate ink reduction amount.\nAlso, US 2005/0063749 fails to disclose a method for setting lightness and saturation of an image to be printed when calibrating data for reducing ink consumption, and it is unclear how data is processed for reducing ink consumption."}
{"text": "Tanks for the storage of liquids have been constructed in a variety of ways from a variety of materials. In one common application, the underground storage of hydrocarbons, such as gasoline and other petroleum products, the tanks have conventionally been fabricated out of steel or fiberglass, most commonly with a single rigid wall. In many applications this construction has proved reasonably satisfactory, with such tanks functioning properly for many years before requiring repair or replacement. However, the increasing age of many of the tanks currently in place is beginning to present serious environmental dangers. Many of the older steel tanks buried underground have rusted and are beginning to leak, thus releasing the petroleum materials into the ground where they may seep into and pollute underground water supplies. While rustproof, some fiberglass tanks have also exhibited leakage, causing the same problems.\nOne of the primary problems with leaking storage tanks has been the difficulty or inability to ascertain when or if such leaks are occurring from a given tank. Because the excavation and removal of such a storage tank, which may contain thousands of gallons of fuel, is an expensive and difficult undertaking, such an operation is difficult to justify unless there is some evidence of actual leakage.\nBecause of the increasing potential danger of leaking storage tanks, particularly in communities that utilize ground water for public consumption, many municipalities have implemented or plan to implement ordinances requiring the use of double wall storage tanks underground and requiring replacement of existing single wall tanks. While the installation of a conventional double wall tank in a new facility entails no great difficulty and a generally manageable increase in cost over a single wall tank, the burden of complying with such ordinances by replacing existing sound, single wall tanks with double wall tanks can be heavy. This burden has prompted the search for methods of fabricating relatively inexpensive double wall tanks. This burden has also given impetus to the search for a method of remanufacturing existing single wall tanks into a double wall assembly with means for detecting the presence of any leaks into the space between the two walls."}
{"text": "The present invention relates to variable displacement vane pumps and, in particular, to pumps having rotors carrying radially slidable vanes.\nVane pumps are well known. In a very simple form, a pumping rotor carries vanes around a casing having a cylindrical volume which is eccentric to the center of rotation of the rotor. This eccentricity results in the volume between vanes changing cyclically. With ports in the casing positioned appropriately, the changing volume between vanes can cause pumping. In a known rotary pump, the extent of eccentricity between rotor and casing can be changed to alter the pumping displacement. While a very successful type of pump, it develops a high radial load on the rotor, requiring massive bearings whose high weight is unsuitable for aircraft and other applications. A balanced vane pump, however, is symmetrical about any diameter and since it applies no radial load to its journal bearings, it is much lighter.\nAnother known pump employs a rotor having along its periphery axial slots carrying vanes which are axially reciprocatable within the slots. A coaxial cam in the shape of a truncated cylinder lies alongside these vanes and directs their axial reciprocation at the rate of one cycle per revolution of the rotor. While this arrangement can be set to change the phasing of the vanes and the associated volume between them, this phasing change has some disadvantages. When the phasing is changed to reduce pump displacement, the fluid being pumped is pressurized and depressurized non-productively, thereby increasing the load on the various moving parts and pump bearings.\nAnother known fluid motor or pump has a duplex construction. This construction includes two cams on opposite sides of a blade carrier, forming two working chambers. The blade carrier has two circumferential series of separate blades. However, this known device does not provide for relative rotation between the different cams to change the displacement of a pump. Thus this design suffers the unnecessary loading mentioned previously.\nIt is also known to alter the displacement of a vane pump by distorting the surface of its cam. A disadvantage with such distortions is that the peak acceleration and thus the forces applied to the vane varies significantly as the cam is distorted. Thus vane wear increases as the cam is distorted to non-ideal configurations.\nAccordingly, there is a need for a variable displacement, preferably balanced, vane pump which has a cam surface that does not create undue acceleration and stress. This pump ought to be able to vary the displacement and reduce the input power when the flow of the pump is reduced. Furthermore the pump ought to be simple, efficient and reliable."}
{"text": "Use of radio frequency bands of the electromagnetic spectrum is regulated by governments in most countries, by allocating specific frequency bands to particular types of uses, such as licensed bands for commercial radio and television broadcasting, cellular telephony, mobile networks such as CDMA2000, WCDMA, HSPA, LTE, and IMT, maritime radio, police, fire, and public safety radio, GPS, radio astronomy, earth stations for satellite communications, and many other uses. Governments also allocate unlicensed bands, for example, for Wireless Regional Area Network (WRAN) broadband access for rural areas and wireless local area networks (WLAN) and wireless personal area networks (WPAN), such as the industrial, scientific, and medical (ISM) band.\nIn the United States, the Federal Communications Commission (FCC) regulates use of the radio spectrum, including radio and television broadcasting. Frequencies are allocated according to a bandplan in which guard bands are assigned between the allocated radio bands to avoid interference between adjacent signals. There are also unassigned frequency bands in the spectrum that either have never been used or have become free as a result of changes in technology. Unassigned or un-used frequencies also appear locally inside the frequency bands, which are otherwise allocated in other locations. The unassigned frequency bands and guard bands are referred to as white spaces.\nTV white space may be broadly defined as broadcast television spectrum that is unused by licensed services. There are at least two categories of TV white space: [1] Dedicated TV white space is a portion of the spectrum that the FCC has reallocated to unlicensed use from previously analog broadcast usage, and [2] Locally unused spectrum by licensed TV broadcasters in a geographic area.\n[1] Dedicated TV white space: In the United States, the FCC has dedicated approximately 400 MHz of white spaces for unlicensed use that became unused after a federally mandated transformation of analog TV broadcasting to digital TV broadcasting. However, the FCC has prohibited unlicensed use of white spaces from interfering with existing licensed uses, including digital TV stations, low power TV stations, cable TV headends, and sites where low power wireless microphones are used. Various proposals have been made for unlicensed use of the white spaces left by the termination of analog TV, for example rural broadband deployment, auxiliary public safety communications, educational and enterprise video conferencing, personal consumer applications, mesh networks, security applications, municipal broadband access, enhanced local coverage and communications, fixed backhaul, and sensor aggregation for smart grid meter reading.\n[2] Locally unused spectrum by licensed TV broadcasters: The FCC has adopted rules to allow unlicensed radio transmitters to operate in the broadcast television spectrum at locations where that spectrum is not being used by licensed broadcasters. The FCC proposes two mechanisms to enable the unlicensed transmitter to discover the available channels: geo-location and database based approach, and spectrum sensing. The use of one of the mechanisms is required for the unlicensed transmitter. The FCC proposed the use of geolocation to establish the location of the unlicensed transmitter and a database of TV bands use by licensed broadcasters organized by their geographic coverage areas, to enable the unlicensed transmitter to know where local TV band white spaces may be available. The FCC proposed the use of spectrum sensors in the unlicensed transmitter to detect the presence of the incumbent, primary TV broadcaster's signal in the local TV band to enable the unlicensed transmitter to immediately relinquish using the band. A primary user in such a local TV band would be an incumbent TV broadcaster licensed to operate in that band, but in those geographic areas where there are no licensed incumbent TV broadcasters in operation, other unlicensed secondary users may make use of that band. There may be also other incumbent users in the TV band, which the secondary users should avoid, such as program making and special events (PMSE) systems.\nIn addition to the United States, other countries are also considering to enable unlicensed, secondary operation in TV band white spaces. The requirements may slightly differ in different countries, e.g. in the United States the maximum transmit power for unlicensed device is defined based on the device type, whereas in Europe location specific maximum transmission power has been considered. In that case the maximum allowed transmission power for an unlicensed device would depend on the device geo-location, i.e. the distance from the primary users. The device characteristics, such as emission mask/ACLR (adjacent channel leakage ratio) may affect the maximum allowed transmission power.\nOther RF spectrum white spaces may be defined as RF spectrum that is locally unused in certain geographic areas, such as for example frequency allocations from maritime radio in landlocked areas remote from the sea. A primary user in such a maritime radio band would be a maritime radio licensed to operate in that band, but in those geographic areas where there are no licensed maritime radios in operation, other unlicensed secondary users may make use of that band. Similarly, locally unused spectrum white spaces may be present in certain geographic locations, such as the frequency allocations from 2.025 GHz to 2.110 GHz for earth stations to transmit to communications satellites, in areas remote from such earth stations. A primary user in such a satellite earth station radio band would be a satellite earth station licensed to operate in that band, but in those geographic areas where there are no satellite earth stations in operation, other unlicensed secondary users may make use of that band. Further, other schemes of secondary use of spectrum, other than unlicensed schemes may exist, such as licensing, regulator defined policies, cognitive principles, or authorized shared access."}
{"text": "Pursuant to 37 CFR 1.71(e), this patent document contains material which is subject to copyright protection and the owner of this patent document reserves all copyright rights whatsoever.\nIn the late 1980's NIRS was determined to be useable to noninvasively assess skeletal muscle oxygenation (SmO2) in patients with heart failure by comparing light absorption at 760 nm and 800 nm as indicia of hemoglobin-myoglobin oxygenation. Wilson, J. R. et al.; Noninvasive Detection Of Skeletal Muscle Underperfusion With Near-Infrared Spectroscopy In Patients With Heart Failure; Circulation, 80(6), Pages 1668-74 (1989). Since then, NIRS has been studied for use in measuring oxygenation concentration as well as other physiological variables and analytes in various organs and tissues of the body. In contrast to typical pulse oximetry, which generally measures oxygen in the flowing blood, NIRS can be used to measure whether or not enough oxygen is being delivered to meet the metabolic demand of a particular organ or tissue (e.g., skeletal muscle, heart, brain, etc.).\nNIRS has been studied for use in monitoring the oxygen content of certain body tissues during and following cardiopulmonary resuscitation (CPR). In one such study, regional cerebrovascular oxygen saturation (rSO2) was monitored by placing an infrared light-emitting probe on the patient's forehead after arrival in the hospital emergency department. Patients who survived for one week had significantly higher median rSO2 on arrival than nonsurvivors. Also, patients who arrived while undergoing CPR without spontaneous circulation had lower median rSO2 than patients who arrived after restoration of spontaneous circulation (ROSC). Patients with ROSC who went on to survive for one week had a higher rSO2 on arrival than patients with ROSC who did not survive for one week. These investigators concluded that low rSO2 after cardiac arrest was associated with a higher mortality and that non-invasive monitoring of cerebrovascular oxygen saturation by NIRS could potentially be useful in prognosticating outcomes for patients following cardiac arrest. Mullner, M., et al., Near Infrared Spectroscopy During And After Cardiac Arrest—Preliminary Results; Clinical Intensive Care, Vol. 6, No. 3, Pages 107-11 (1995).\nIn a more recent study, patients who had experienced out-of hospital cardiac arrest followed by ROSC were monitored by an NIRS StO2 monitor and by an end-tidal carbon dioxide (ETCO2) monitor. ETCO2 had previously been established and an indicator of ROSC or rearrest. Downward trends in StO2 were observed prior to each rearrest and rapid increases in StO2 were noted after ROSC. The StO2 data showed less variance than the ETCO2 data in the periarrest period. The investigators concluded that a decline in StO2 level may correlate with rearrest and, thus. may be useful as a predictor of rearrest in post-cardiac arrest patients. A rapid increase in StO2 was also seen upon ROSC and may be a better method of identifying ROSC during CPR than pauses for pulse checks or ETCO2 monitoring. Frisch, A., et al.; Potential Utility of Near-Infrared Spectroscopy in Out-of-Hospital Cardiac Arrest: An Illustrative Case Series; Prehospital Emergency Care, Vol. 16, No. 4: Pages 564-570 (2012).\nAdditionally, investigators have explored the use of a subcutaneously implanted NIRS device in combination with an Implanted Cardioverter Defibrillator (ICD). In this study, NIRS oximetric measurements were used, in combination with electrical monitoring by the ICD, to distinguish between the onset of a ventricular arrhythmia requiring defibrillation and mere electromagnetic interference or artifacts resulting from erroneous double counting of the electrocardiographic T-wave as an R-wave, ICD lead failure, or other electrocardiographic aberrancies. Bhunia, S. K. et al., Implanted Near-Infrared Spectroscopy For Cardiac Monitoring; Proc. SPIE 7896, Optical Tomography and Spectroscopy of Tissue IX, 789632 (2011). [http://proceedings.spiedigitallibrary.org/proceeding.aspx?articleid=733147]\nThe prior art has included a number of NIRS devices that are positionable at various locations on the patients skin, or subcutaneously below the skin, to measure physiological properties or concentrations of analytes (e.g., pH, temperature, oxygen tension, oxygen saturation, partial pressure of oxygen, partial pressure of carbon dioxide, hemoglobin concentration, water concentration, hematocrit, glucose concentration, presence of biomarkers, etc.) in underlying organs or tissues. Some but not necessarily all examples of such devices are described in U.S. Pat. No. 5,931,779 (Arakaki, et al.); U.S. Pat. No. 6,212,424 (Robinson); U.S. Pat. No. 6,990,364 (Ruchti et al.); U.S. Pat. No. 7,245,373 (SoIler, et al.); U.S. Pat. No. 7,613,489 (Myers); U.S. Pat. No. 7,647,092 (Motz et al.); U.S. Pat. No. 8,277,385 (Berka et al.); U.S. Pat. No. 8,346,329 (Xu et al.); U.S. Pat. No. 8,406,838 (Kato) and U.S. Pat. No. 8,649,849 (Liu et al.) as well as United States Patent Application Publication Nos. 2014/0135647 (Wolf II); 2014/0024904 (Takinami); 2013/0225955 (Schenkman, et al.) and 2011/0184683 (SoIler et al.), the entire disclosure of each such patent and patent application being expressly incorporated herein by reference. Also, examples of such devices are currently marketed as CareGuide™ Oximeters (Reflectance Medical, Inc., Westborough, Mass.); INVOS™ Somatic/Cerebral Oximetry Monitors (Covidien Respiratory and Monitoring Solutions, Boulder, Colo.); Reveal LINQ™ Insertable Cardiac Monitoring Systems (Medtronic Corporation, Minneapolis, Minn.); FORE-SIGHT ELITE® Cerebral Oxygen Monitors (CAS Medical Systems, Inc., Branford, Conn.) and EQUANOX™ Cerebral/Somatic Tissue Oximetry Devices (Nonin Medical, Inc., Plymouth, Minn.). Some if not all of these NIRS devices utilize specialized apparatus and/or signal processing techniques (e.g, “background subtraction”) to minimize or eliminate spectral effects from skin, bone or other intervening tissue that resides between the location of the NIRS device and the organ or tissue of interest.\nIn the past, certain devices have been positioned within the esophagus adjacent to the heart to monitor or image the heart from a vantage point that has minimal intervening tissue between the device and the heart. For example, endoesophageal stethoscopes and pulse monitoring probes have been advanced into the esophagus and used to monitor a patient's heartbeat, examples of which are described in U.S. Pat. No. 4,409,986 (Apple et al.); U.S. Pat. No. 4,331,156 (Apple, et al.). Also, ultrasound probes have been inserted into the esophagus and used for transesophageal echocardiography, examples of which are described in U.S. Pat. No. 8,641,627 (Roth et al.); U.S. Pat. No. 8,172,758 (Harhen); U.S. Pat. No. 6,884,220 (Aviv et al.) and U.S. Pat. No. 6,471,653 (Jordfald, et al.).\nAdditionally, United States Patent Application Publication No. 2013/0231573 (Zeng et al.) describes the insertion of a near-infrared spectroscopy probe through the working channel of a bronchoscope and the use of such probe to for endobroncheal Raman spectroscopic analysis of lung cancer tissue.\nGiven that NIRS and other forms of optical spectroscopy used for measuring the physiologic status of living tissue (broadly termed “Physiologic Spectroscopy” (PS)) are a potentially valuable tool for monitoring cardiac tissue or a subject's blood at specific locations in emergency and critical care situations, it is desirable to develop new optical spectroscopy monitoring devices and methods which are useable for obtaining optical spectrographic measurements from cardiac tissue or from blood located within the chambers of the heart or great vessels (e.g., pulmonary artery, aorta, etc.), and other tissues within the vicinity of the esophagus, the trachea or the main bronchi."}
{"text": "1. Field of the Invention\nThis invention relates to apparatus for pumping fluids intravenously or intra arterially to patients, and more particularly, it relates to the control systems for such fluid infusion pumping apparatus which regulate the flow rates under which such pumping is accomplished.\n2. Description of the Prior Art\nIn recent years there has been a considerably increased use of positive displacement fluid infusion pumping devices for delivering fluids intravenously or intra arterially to patients in hospitals or other patient care locations. These have to a large extent replaced the time honored gravity flow control systems primarily due to their much greater accuracy in determining delivery rates and dosages, their relative sophistication in permitting a flexible and controlled feed from multiple liquid sources, and particularly their ability to control with precision the amount of dangerous drugs delivered to a patient over a given period of time.\nIn a typical positive displacement fluid infusion pump system, stepping motors are used to gradually deliver the I.V. fluids by displacing an amount of liquid trapped within a chamber in very small increments. In this way, shock to the nervous system of a patient is minimized and a relatively uniform flow of fluid is obtained. Various prior United States patents disclosing typical fluid infusion pump control systems include U.S. Pat. Nos. 4,474,309 to Edward G. Solomon; 4,443,216 to Anthony G. Chappell; 4,396,385 to Timothy G. Kelly et al.; 4,392,849 to John H. Petre et al.; 4,308,866 to Roger W. Jelliffe et al.; 4,280,494 to Robert J. Cosgrove, Jr. et al.; 4,217,993 to Thurman S. Jess et al.; 4,126,132 to Peer M. Portner et al.; 4,037,598 to Heinz W. Georgi; and 3,985,133 to John Arthur Jenkins et al.\nThe United States patent to Solomon, U.S. Pat. No. 4,474,309, discloses a typical fluid infusion pump system wherein a small disposable cassette having a pumping chamber with a flexible diaphragm forming one exterior surface thereof is adapted to be placed in a driver or controller mechanism, and with the cassette including a programmable plunger adapted to interact with the diaphragm to pump fluid from the chamber. The plunger is driven by a stepping motor through a series of short, liquid-displacing steps and then retracted quickly once it has reached its maximum displacement in the chamber so as to assure as uniform a flow of liquid as possible. The inlet and outlet valves to the pumping chamber may be actuated from the driver (as is the case with the inlet valve in the Solomon patent disclosure) or they may be comprised of check valves controlled by the differential pressure provided by the pumping plunger (as in the case with the outlet valve in the Solomon disclosure patent). As is conventional with pumping systems of this type, the user sets a rate at which the fluid is to be delivered as well as the total dosage which is to be delivered by the apparatus. The system will then operate by providing the fluid at the predescribed rate in a generally uniform flow and will automatically shut down when the accumulated running total of liquid delivered equals the preset total dosage to be delivered.\nWhile careful control of the pumping mechanism and limiting the delivery to short strokes spaced over time insures a more uniform flow rate, it has been found that physiological shocks may still occur in particular patients, especially when the delivery of I.V. solutions is started and also when it abruptly stops."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method for converting an interlaced image into a progressive image, and more particularly, to an apparatus and method for converting an interlaced image into a progressive image by interpolating the interlaced image considering characteristics of a region and the orientation of an edge to which a pixel to be interpolated belongs.\n2. Description of the Related Art\nThere are various methods of converting an interlaced image into a progressive image. In particular, a blend method and an edge dependent interpolation method are mainly used.\nIn the blend method, interpolation is performed on an object pixel, using the values of pixels on preceding and following interlaced scan lines corresponding to the object pixel and the values of upper and lower pixels above and below the object pixel, and the result of interpolation is used as the value of the object pixel. The blend method is preferred since an edge of an image can be smoothly and naturally represented. However, when an interlaced image with much motion is converted into a progressive image using the blend method, a blur and an afterimage are likely to occur, and diagonal aliasing may occur at an edge of the image.\nIn the edge dependent interpolation method, the value of an object pixel, which is to be interpolated, is obtained using the values of pixels above and below the object pixel and the values of pixels to the left and right of the pixels above and below the object pixel.\nReferring to FIG. 1, a value obtained by progressively scanning an object pixel X is computed using the value Dlu of a pixel to the upper left of the object pixel, the value Drd of a pixel to the lower right of the object pixel, the value Dru of a pixel to the upper right of the object pixel, and the value Dld of a pixel to the lower left of the object pixel, and the values Vu and Vd of pixels above and below the object pixel.\nMore specifically, in the edge dependent interpolation method, the value obtained by progressively scanning an object pixel X is obtained by computing the differences between the values Vu and Vd, between the values Dlu and Drd, and between the values Dru and Dld, and then interpolating the object pixel X using pixels disposed in a direction having the smallest value among these differences. This method is advantageous in that slanted edges of an image can be naturally represented since image conversion is made in consideration of pixels disposed diagonally from the object pixel, thus preventing aliasing from occurring in the image.\nHowever, the edge dependent interpolation method is disadvantageous in that the edge orientation in an image is likely to be erroneously determined and use of only one field of an interlaced image during interpolation may cause a large difference between chromaticity values of the interlaced image and a progressive image. Also, color blurring may occur during conversion of an image having an object moving rapidly into a progressive image. Accordingly, the edge dependent interpolation method requires determining whether a pixel is to be interpolated in a diagonal direction or a vertical direction.\nConventionally, the edge orientation in an image is determined by computing absolute values |x(k)−y(−k)| of the differences between values of pixels along an upper scan line and pixels of a lower scan line (k=−N, . . . , 0, . . . , N), and determining that there is a diagonal line in the direction in which the smallest absolute value is obtained. Such a method is disclosed in U.S. Pat. Nos. 5,532,751 and 6,421,090. However, this method requires complex computation to obtain the absolute values of the differences between the values of possible combinations of pixels disposed diagonally from an object pixel. When computing absolute values in only several diagonal directions to reduce computational complexity, defects may occur in an image. Also, motion estimation needed for precise image conversion requires a lot of computation and large memory bandwidth, thus making it difficult to convert an interlaced image into a progressive image in real time."}
{"text": "Early stage prostate cancer diagnosis and staging challenges patients and clinicians, as it has proven difficult to reliably distinguish indolent and incidental disease from progressive, life-threatening disease. Standard treatment for prostate cancer generally has been directed at the entire organ or perhaps specific organ regions found generally to have a higher incidence of tumor burden, rather than at point source identification of pathologically confirmed regions of tumor. Traditionally, organ-targeted prostate cancer treatments have been accomplished either by surgical resection of the gland, or by directing multiple beams of radiation into the pelvis to encompass the entire gland in a uniform dose. In the case of prostate cancer, tumor is not partially resected surgically for treatment, as is the case in breast and other cancers. Rather, the entire prostate gland is removed during surgical intervention. Radiation treatments for prostate cancer also treat the gland and glandular areas found to have higher incidence of tumor burden, often having the impact of over-exposing normal radiosensitive tissues to unnecessary treatment, thereby increasing the likelihood of treatment related morbidities. Such poorly directed curative therapies lead to increased cost of care and decreased patient quality of life.\nAlternate methods such as cryotherapy have also been explored, though again generally directed to treat, or freeze, the entire organ. While newer focal cryotherapy methods attempt to more uniquely focus the freezing process to a section of the gland in order to decrease the incidence of treatment related morbidity, appropriate targeting methods continue to rely on general assumptions about the location of areas of high tumor burden.\nAs computers have become better at assisting in treatment planning, it has become standard to use three dimensional (3-D) treatment planning to shape the individual field for radiation therapy to treat, for example, only the prostate itself, commonly utilizing a minimal margin of about 5 to 15 millimeters beyond the edges of the gland. Use of this planning convention allows for higher doses of radiation to be delivered while at the same time sparing the surrounding tissue and limiting morbidity. Clinical trials have confirmed that by using higher doses of radiation in a 3-D conformal approach, higher cure rates are achieved with lower toxicity.\nAn alternative method for achieving a high dose delivered to the entire gland is known as brachytherapy, or the placement of either temporary High Dose Rate (HDR) or permanent Low Dose Rate (LDR) radioactive sources within the prostate. These therapies allow focal radiation therapy targeting to within 10 millimeters of the glandular areas generally identified as at high risk of containing tumor. Target treatment margins for brachytherapy are typically adopted in the range of 1-10 mm. A gross target volume (GTV) of a treatment target may be defined by anatomic image studies, which may then be used to further define a clinical target volume (CTV), typically comprising the GTV plus an adequate margin to account for microscopic disease at the edge of the GTV and allowance for motion of the GTV from patient positioning variation during image study. Rarely are permanent or temporary sources placed more than 10 mm beyond either the CTV or GTV. For brachytherapy (LDR and HDR), the CTV often is equal to the GTV as there is no motion of the organ that does not include the sources, and daily set-up errors can be eliminated. A biological target volume (BTV) typically represents a region defined by a functional study that may be completely within the GTV, or may expand the GTV by showing disease extending beyond the margins defined by the GTV on the anatomic study.\nMore recently, again enabled in large part through improved computer software, a newer external beam radiation therapy technique, referred to as Intensity Modulated Radiation Therapy or IMRT, has become available for treating the entire organ with tighter margins of as little as 4 millimeters. IMRT provides options for targeting small volume (<1 cc) regions within a treatment planning volume (TPV) to focus higher doses than the dose delivered to the entire gland volume, comprised of the CTV, GTV and BTV. This focused IMRT treatment method, as described, is currently utilized in only a minority of select academic settings using functional images acquired with either Magnetic Resonance Spectroscopy Imaging (MRSI) or Single Photon Emission Computerized Tomography (SPECT) images to help define a region within the prostate gland believed to represent occult tumor volumes. These identified areas found to be suspicious for occult tumor on functional imaging represent findings which are indistinguishable with standard anatomic studies such as Computerized Axial Tomography (CAT or CT) scan, Magnetic Resonance Imaging (MRI) used in conjunction with Ultrasound (US) and/or US alone. While the SPECT imaging techniques rely on over expression of a specific protein identified by a radiolabeled monoclonal antibody, the MRSI technique utilizes voxel analysis of tissue composition to detect regions felt more likely to represent cancerous regions. Newer functional studies will certainly be developed in the future using similar technologies, such as Positron Emission Tomography (PET) tracers, Optical Biopsy techniques, or other similar technologies.\nAs discussed in more detail below, the method of the invention will enable the more effective utilization of the foregoing image modalities to localize and treat cancer within a particular target organ with tumor site localization confirmation to histopathological findings in order to more effectively target dose escalation to BTV while sparing surrounding tissues from unnecessary treatment.\nBecause standard image techniques (CAT, MRI, X-ray, US) are unable, in routine clinical use, to visualize specific regions containing tumor within the gland, it remains a significant problem for current therapies, resulting in increased treatment related morbidity and overall cost of care. In routine clinical practice, at the time a male patient presents with either a palpable abnormality on Digital Rectal Exam (DRE) or, more commonly, with an elevated Prostate Specific Antigen (PSA) level on a blood test, a biopsy of the gland is recommended by the physician to determine if a cancer is present. The patient's initial biopsy procedure is typically performed in the office of a urologist. Frequently, this first procedure is completed with the patient in lateral decubitus position utilizing a Trans-Rectal Ultrasound (TRUS) probe that allows biopsy sample to be taken through the TRUS probe. Sextant biopsy is regarded as the standard of reference for nonsurgical tumor localization, although limitations of sextant biopsy are increasingly recognized.\nDuring a standard sextant biopsy procedure, typically six to twelve biopsy tissue samples will be obtained from the prostate gland with each biopsy sample involving an individual needle pass through the rectal wall and into the desired location within the gland for the biopsy. Standard sextant biopsies are directed into both the right and left general regions of the prostate gland, and may further be directed into the right and left base, mid and/or apex regions of the gland at either medial or lateral locations. Recently, it has become more common to have biopsy tissue samples recorded as to the rudimentary region within the gland from which it was obtained (medial/lateral-right/left: base, mid or apex). In addition, pathologists are more frequently in standard practice being requested to record and report the percentage of each biopsy core involved with tumor to help determine if the region of biopsy has minimal disease or bulky tumor deposits.\nWhen a patient presents as highly suspicious for disease (e.g., rising PSA or positive DRE) and disease confirmation cannot be validated by a positive biopsy result, patients more commonly today undergo saturation biopsy procedures whereby 24-36 biopsy tissue samples are taken, making it increasingly difficult and costly to record and track the region within the gland from which the sample was obtained. Most often, only positive biopsy samples are reported for rudimentary location. Most importantly, the precise location of each biopsy sample cannot be determined even with the use of imaging studies following the biopsy procedure because the sextant localization of disease is not synonymous with volumetric localization of tumor. As such, there is no effective means of correlating the specific pathology of the biopsy site to its location in the target organ or tissue in a manner that will effectively facilitate precise localization of positive histopathology to identify tumor volumetric localization for use in treatment planning targeting. Cancer patients, such as prostate, are frequently followed for extended time periods between diagnosis, medical imaging, treatment and post-therapy follow-up. These patients are, therefore, evaluated over time by different physicians (e.g., urologists and radiation oncologists) in a number of settings (e.g., physician office, outpatient hospital imaging, surgical center) and with various imaging and image-guided treatment modalities requiring different patient positioning during imaging (supine) and treatment (lateral decubital) which, collectively, serve to obscure correlation of pathology information with in-vivo image sets.\nThus, there remains a need for the ability to refine the identification and volumetric location of disease as correlated to positive histopathology within the suspect organ or tissue, and to be able to use these refined data sets to direct therapy in a more effective and less harmful manner. It has been estimated that each year in the U.S. over one million biopsies are performed, with as many as 50% of those core samples being reported as negative. In the example of standard core biopsy sampling techniques for the prostate gland collecting 6 to 12 standard core samples per patient, millions of core tissue biopsy samples are evaluated each year. The clinical inefficiency of the biopsy procedure results in a large negative burden to patients for non-productive procedures correlating to increased risk for procedure related morbidities (e.g., fever, infection and bleeding, discomfort, and lost productivity), for the often ineffective procedure. In addition, the pathology results are routinely lost to discreet anatomic localization for use in therapy planning."}
{"text": "The present invention relates to a device for loading signatures and the like to be applied to sheet loading machines for bookbinding use.\nMore specifically the device affords the possibility of progressively reducing, by partially narrowing the feeding line, the weight applied on the signatures being conveyed to the sheet loading assembly, as said signatures are advanced towards a locating conveyor belt.\nAs is known, in a bookbinding process the signatures, that is the already bent sheets to be printed upon, are picked by means of a sheet or signature loading assembly which conveys them to a processing apparatus.\nThe mentioned picking operation, on the other hand, is affected by drawbacks because of the practical difficulties associated with disengaging the signatures from one another.\nThere are known signature loading apparatus which are effective to suitably feed, by means of a conveyor belt, the thereinabove mentioned sheet or signature loading assemblies.\nHowever, these loading apparatus are not able of satisfactorily operating, since they feed the mentioned conveyor belt, usually obliquely arranged, by means of supplying means entraining or pushing the signatures, adjoining one another and in close mutual contact, towards said conveyor belt by causing said signatures to fall, upon tilting, on the conveyor belt itself.\nThus, since a proper signature disengaging step is lacking, said signatures may be unevenly conveyed to the sheet or signature loading assembly with consequent interruptions in the operation of the system."}
{"text": "The present invention relates to a new and distinct cultivar of Osteospermum plant, botanically known as Osteospermum ecklonis, and hereinafter referred to by the name ‘Fidostsepia’.\nThe new Osteospermum plant is a product of a planned breeding program conducted by the Inventor in De Lier, The Netherlands. The objective of the program is to create and develop new compact Osteospermum plants that are freely branching with numerous inflorescences with unique and attractive ray and disc floret coloration.\nThe new Osteospermum plant originated from a cross-pollination by the Inventor in May, 2008 of a proprietary selection of Osteospermum ecklonis identified as code number O 45711, not patented, as the female, or seed, parent with a proprietary selection of Osteospermum ecklonis identified as code number O 67242, not patented, as the male, or pollen, parent. The new Osteospermum plant was discovered and selected by the Inventor as a single flowering plant within the progeny of the stated cross-pollination in a controlled greenhouse environment in De Lier, The Netherlands in October, 2008.\nAsexual reproduction of the new Osteospermum plant by terminal cuttings in a controlled greenhouse environment in De Lier, The Netherlands since November, 2008 has shown that the unique features of this new Osteospermum plant are stable and reproduced true to type in successive generations."}
{"text": "1. Field of the Invention\nThis invention relates to a support framework for a flat panel display as represented by a liquid crystal panel which is an object to be tested, or it relates to a support framework for a testing probe block which is subjected to pressure contact with electrode pads which are arranged on a peripheral edge portion of the flat panel display.\n2. Related Art\nThe flat panel display support framework and testing probe block support framework are fabricated of a rectangular framework. As shown in FIGS. 6(A) and 6(B), a front rectangular framework 1, on which testing probe blocks 8 are supported, is disposed on a front side of a rear rectangular framework 1, on which a flat panel display 7 as represented by a liquid crystal is supported. The rear rectangular framework 1 is reciprocally movably disposed on the front rectangular framework 1. The arrangement is such that when the rear rectangular framework 1 is advanced, testing probes 9 of the testing probe blocks 8 are pressure contacted with corresponding electrode pads 10 which are arranged on a peripheral edge portion of the flat panel display 7. When the rear rectangular framework 1 is retracted, replacement of the flat panel display 7 is performed.\nIt is ordinary practice in the prior art that a rear and a front framework are prepared for each of the flat panel displays which have difference sizes, and a testing probe block is supported thereon.\nThis is inconvenient, indeed, because it is necessary to prepare a test device having proper front and rear frameworks for each of the flat panel displays which have different sizes. Otherwise, it is necessary to design only the front and rear frameworks replaceable in accordance with the sizes of the flat panel displays. This naturally results in not only increased expenses in installation of the device but also inefficiency of testing operation carried out by a worker(s) who is required to replace heavy frameworks.\nIt is, therefore, a general object of the present invention to provide a flat panel display support framework or a probe block support framework which can amicably solve the above-mentioned inconveniences which the conventional devices had.\nA specific object of the present invention is to provide a flat panel display support framework or a probe block support framework which can appropriately cope with flat panel displays having different sizes in a single test device.\nIn order to achieve the above objects, there is essentially provided in a first embodiment of the present invention a rectangular support framework for a flat panel display as an object to be tested or a testing probe block. The rectangular support framework comprises an upper cross frame plate which is capable of moving vertically alone with respect to a right vertical frame plate and which is capable of moving vertically in association with vertical movement of a left vertical frame plate, a lower cross frame plate which is capable of moving vertically alone with respect to the left vertical frame plate and which is capable of moving vertically in association with vertical movement of the right vertical frame plate, the left vertical frame plate which is capable of moving laterally alone with respect to the upper cross frame plate and which is capable of moving laterally in association with lateral movement of the lower cross frame plate, and the right vertical frame plate which is capable of moving laterally alone with respect to the lower cross frame plate and which is capable of moving laterally in association with the lateral movement of the upper cross frame plate. The vertical movement and lateral movement of the upper cross frame plates and left and right vertical frame plates enlarge or reduce a window which is defined by the upper and lower cross frame plates and left and right vertical frame plates. The flat panel display or the probe block which is subjected to pressure contact with electrode pads arranged on a peripheral edge portion of the flat panel display is supported on the upper and lower cross frame plates and left and right vertical frame plates which define the enlarged or reduced window.\nIn the first embodiment, the upper cross frame plate and the right vertical frame plate are connected such that the upper cross frame plate is capable of moving vertically alone and the upper cross frame plate and right vertical frame plate are capable of moving laterally in association with each other, the left vertical frame plate and the upper cross frame plate are connected such that the left vertical frame plate is capable of moving laterally alone and the left vertical frame plate and upper cross frame plate are capable of moving vertically in association with each other, the lower cross frame plate and the left vertical frame plate are connected such that the lower cross frame plate is capable of moving vertically alone and the lower cross frame plate and left vertical frame plate are capable of moving laterally in association with each other, and the right vertical frame plate and the lower cross are connected such that the right vertical frame plate is capable of moving laterally alone and the right vertical frame plate and lower cross frame plate are capable of moving vertically in association with each other.\nIn a second embodiment of the present invention, there is provided a rectangular support framework for a flat panel display as an object to be tested or a testing probe block. The rectangular support framework comprises an upper cross frame plate which is capable of moving vertically alone with respect to a left vertical frame plate and which is capable of moving vertically in association with vertical movement of a right vertical frame plate, a lower cross frame plate which is capable of moving vertically alone with respect to the right vertical frame plate and which is capable of moving vertically in association with vertical movement of the left vertical frame plate, the left vertical frame plate which is capable of moving laterally alone with respect to the lower cross frame plate and which is capable of moving laterally in association with lateral movement of the upper cross frame plate, and the right vertical frame plate which is capable of moving laterally alone with respect to the upper cross frame plate and which is capable of moving laterally in association with lateral movement of the lower cross frame. The vertical movement and lateral movement of the upper and lower cross frame plates and left and right vertical frame plates enlarge or reduce a window which is defined by the upper and lower cross frame plates and left and right vertical frame plates. The flat panel display or the probe block which is subjected to pressure contact with electrode pads arranged on a peripheral edge portion of the flat panel display is supported on the upper and lower cross frame plates and left and right vertical frame plates which define the enlarged or reduced window.\nIn the second embodiment the upper cross frame plate and the left vertical frame plate are connected such that the upper cross frame plate is capable of moving vertically alone and the upper cross frame plate and left vertical frame plate are capable of moving laterally in association with each other, the right vertical frame plate and the upper cross frame plate are connected such that the right vertical frame plate is capable of moving laterally alone and the right vertical frame plate and upper cross frame plate are capable of moving vertically in association with each other, the lower cross frame plate and the right vertical frame plate are connected such that the lower cross frame plate is capable of moving vertically alone and the lower cross frame plate and right vertical frame plate are capable of moving laterally in association with each other, and the left vertical frame plate and the lower cross frame plate are connected such that the left vertical frame plate is capable of moving laterally alone and the left vertical frame plate and lower cross frame plate are capable of moving vertically in association with each other.\nThe several features and advantages of the present invention will become more apparent to those skilled in the art by reference to the detailed description which follows taken in conjunction with the several figures of the drawing."}
{"text": "In order to join two tubes when the tubes are not coaxial it is necessary to either use a special coupling or to bend the tubes. Couplings are standard that allow joints to be made at different standardized angles, typically 30.degree., 45.degree., 60.degree., 90.degree., and 180.degree., but when intermediate angles must be made the installer is normally forced to bend the tubes or make up a special coupling.\nIn German patent document 3,914,420 filed May 1, 1989 by A. W. Sommer the coupling ball is hollow and is formed with an array of holes. Spreadable pins are inserted into the tube ends and then into the appropriate holes and the unused holes are capped. This arrangement provides a wide range of positions, but still does not accommodate intermediate angular settings. In addition the capped holes are obvious and are often considered ugly.\nGerman patent document 4,336,401 describes a system where two part-spherical parts are joined together at a diametral plane and each have a mounting pin that is fitted into a respective tube end. When the two tubes are coaxial the diametral plane forms an acute angle to the axis. Thus when an angle needs to be formed, the two coupling parts can be rotated relative to each other so their pins are not coaxial but instead their axes meet at a center of the sphere formed by the two parts. Such a system does indeed allow many angles to be accommodated, but has the disadvantage that the joint in the coupling ball is visible. Furthermore it is fairly difficult to join the two parts together in a manner where the attachment means is not visible."}
{"text": "A fitness application running on mobile device that includes a global navigation satellite system (GNSS) receiver can use GNSS speed and estimated stride length to determine a distance traveled by a user carrying or wearing the mobile device. In some cases, GNSS speed estimates can be lower than the actual speed of the user due to the location of the mobile device on the user's body. For example, if the mobile device is located on the small of the user's back (e.g., stored in a running belt), line-of-sight GNSS signals are obscured by the user's body. Because line-of-sight GNSS signals are blocked by the user's body, the GNSS receiver tracks GNSS signal reflections. Tracking GNSS signal reflections results in lower GNSS speed estimates, which results in less accurate distance traveled calculations."}
{"text": "The present invention relates to α-synuclein antibodies and their use in the prevention or treatment of disease, in particular alpha-synucleinopathies, and more particularly Parkinson's disease (PD).\nAlpha-synucleinopathies, also known as Lewy body diseases (LBDs), are a family of neurodegenerative diseases that all have at their core alpha-synuclein as the key pathological hallmark (Jellinger, Mov Disord (2003), 18 Suppl 6: S2-12, and Spillantini and Goedert, Ann N Y Acad Sci (2000), 920: 16-27; both of which are incorporated herein by reference). Alpha-synucleinopathies include Parkinson's disease (PD), dementia with Lewy bodies (DLB) and multiple system atrophy (MSA).\nPD is a slowly progressive age-related movement disorder affecting greater than 1% of people over 65 years old. PD is the second most common neurodegenerative condition after Alzheimer's disease.\nA defining hallmark pathology of alpha-synucleinopathies are Lewy bodies and Lewy neurites, which are insoluble inclusions of aggregated proteins found inside neurons of the brain revealed upon post-mortem histopathological examination.\nThe presence of Lewy pathology and neuronal loss in non-motor brain regions such as the basal forebrain, mesopontine system, amygdala, neocortex, dorsal motor nucleus of the vagus nerve, olfactory bulbs, locus coeruleus, and the brainstem, may cause cognitive deficits and dementia, hyposmia, sleep disturbances including rapid eye movement sleep behaviour disorder (RBD), mood disorders including depression and anxiety, autonomic dysfunction including cardiovascular and gastrointestinal problems such as constipation, and fatigue and somnolence. Some of these non-motor symptoms appear to characterise the premotor or prodromal phase of the disease (Kalia et al. Lancet (2015), 386(9996): 896-912; incorporated herein by reference).\nThe presence of Lewy pathology and neuronal loss in motor brain regions, including most notably the death of dopaminergic neurons in the substantia nigra, may cause resting tremor, rigidity, bradykinesia and postural instability (Spillantini and Goedert, Ann N Y Acad Sci (2000), 920: 16-27; incorporated herein by reference).\nAlpha-synuclein (also called “α-synuclein” or “α-syn”) protein is the major structural component of Lewy bodies and Lewy neurites. Alpha-synuclein is a small acidic protein made up of 140 amino acids (14 kDa). Human natural wild-type alpha-synuclein has the amino acid sequence SEQ ID NO: 1 as described under UniProtKB accession number P37840. Unless otherwise apparent from the context, reference to alpha-synuclein or its fragments includes the natural human wild-type amino acid sequence indicated above, and human allelic variants thereof, in particular those associated with Lewy body disease (e.g., E46K, A30P, H50Q, G51D and A53T, where the first letter indicates the amino acid in SEQ ID NO: 1, the number is the codon position in SEQ ID NO: 1, and the second letter is the amino acid in the allelic variant). Such variants can optionally be present individually or in any combination. The induced mutations E83Q, A90V, A76T, which enhance alpha-synuclein aggregation, can also be present individually or in combination with each other and/or with human allelic variants E46K, A30P, H50Q, G51D and A53T. At the structural level, alpha-synuclein contains three distinct regions: an amphipathic N-terminal alpha-helix domain that has lipid and membrane binding properties (residues 1-60), a central hydrophobic amyloid-binding domain that encodes the non-amyloid-beta component (NAC) of plaques (residues 61-95), and an acidic proline-rich C-terminal tail (residues 96-140). Residues 71-82 of the NAC domain are believed to be key to the aggregation/fibrillation properties of alpha-synuclein by enabling the protein to switch from a random coil structure to a beta-sheet structure (Bisaglia et al. FASEB J (2009), 23(2): 329-40; incorporated herein by reference). Although the C-terminal domain is free of significant secondary structure it contains a key phosphorylation site at residue Ser129 and a number of tyrosine residues that are nitrated in cytosolic alpha-synuclein inclusions. N-terminal and C-terminal truncated forms of alpha-synuclein also exist. Post-translational modifications to the protein can affect alpha-synuclein aggregation and toxicity (Oueslati et al. Prog Brain Res (2010), 183: 115-45, incorporated herein by reference).\nAlpha-synuclein is abundant in the central nervous system (CNS)/brain where it is found both intracellularly in neurons and glia and also extracellularly in cerebrospinal fluid (CSF) (Mollenhauer el al. J Neural Transm (2012), 119(7): 739-46; incorporated herein by reference) and the interstitial fluid (ISF) that bathes and surrounds the cells of the brain (Emmanouilidou et al. PLoS One (2011), 6(7): e22225; incorporated herein by reference). Alpha-synuclein is a synaptic protein predominantly expressed in neurons of the neocortex, hippocampus, substantia nigra, thalamus, and cerebellum (Iwai el al. Neuron (1995), 14: 467-475; incorporated herein by reference). Under physiological conditions, it is located in neuronal synaptic terminals and is specifically up-regulated at presynaptic terminals during acquisition-related synaptic rearrangement (Fortin et al. J Neurosci (2005), 25: 10913-10921; incorporated herein by reference).\nIn-vitro studies have shown that alpha-synuclein monomers may form the starting point for the aggregation process. The monomer can aggregate into a variety of small oligomeric species that are then stabilised by beta-sheet interactions, going on to form protofibrils which can polymerise into insoluble fibrillary structures reminiscent of those identified in Lewy bodies (Cremades et al. Cell (2012), 149(5): 1048-59; incorporated herein by reference).\nUnder pathological conditions, aberrant alpha-synuclein aggregation may be key to the pathological changes seen in alpha-synucleinopathies (Lashuel et al. Nature (2002), 418: 291, and Tsigelny et al. FEBS Journal (2007), 274: 1862-1877; both of which are incorporated herein by reference). In vitro and in vivo studies have shown that the neurotoxic effects of alpha-synuclein appear to be elicited by small soluble oligomeric conformers or protofibrils (Winner et al. Proc Natl Acad Sci USA (2011), 108(10): 4194-9; and Danzer el al. J Neurosci (2007), 27(34): 9220-32; both of which are incorporated herein by reference). While fibrillar aggregates of alpha-synuclein are characteristic of PD, oligomeric forms of alpha-synuclein are the toxic species (Danzer et al. J Neurosci (2007), 27(34): 9220-32; Lashuel et al. Nature (2002), 418; 291, and Winner et al. Proc Natl Acad Sci USA (2011), 108; 4194-4199; each of which are incorporated herein by reference).\nAlpha-synuclein oligomers can be released to the extracellular environment and taken up by neighboring cells in a “propagation” mechanism (Angot and Brundin, Parkinsonism Relat Disord (2009), 15 Suppl 3: S143-147; Desplats el al. Proc Natl Acad Sci USA (2009), 106: 13010-13015; and Lee el al. J Biol Chem (2010), 285: 9262-9272; each of which are incorporated herein by reference). Aggregates of alpha-synuclein can propagate misfolding through a prion-like spreading mechanism (Lee et al. Nat Rev Neurol (2010), 6: 702-706; Luk et al. J Exp Med (2012), 209(5): 975-86; and Luk et al. Science (2012), 338(6109): 949-53; each of which are incorporated herein by reference). Alpha-synuclein can therefore induce neurodegeneration by either oligomer toxicity or propagation and prion-like spreading.\nIt is now well established and accepted that cells including neurons can secrete various forms of alpha-synuclein (monomers, oligomers, aggregates) under normal conditions and also under conditions of cellular stress, the secretion of monomeric and aggregated forms of alpha-synuclein is elevated under conditions of cellular stress, and through this release of alpha-synuclein into the extracellular milieu, pathological transmissible forms of alpha-synuclein may be propagated between neurons (Recasens and Dehay, Front Neuroanat (2014), 8: 159; incorporated herein by reference).\nThe effects of alpha-synuclein in PD may extend beyond the immediate damage to vulnerable neuronal cells. Like most neurodegenerative diseases there is also a pro-inflammatory cellular response observed (Lee et al. J Biol Chem (2010), 285: 9262-9272; incorporated herein by reference). Circulating alpha-synuclein and/or activated astrocytes can activate microglia, leading to increased generation of reactive oxygen species, nitric oxide and cytokine production, and further exacerbating neurodegeneration (Lee et al. J Biol Chem (2010), 285: 9262-9272; incorporated herein by reference).\nA variety of different experimental models have demonstrated cell-to-cell transmission of alpha-synuclein in cultured cells, or in vivo spreading and propagation of alpha-synuclein pathologies. Lewy body pathology has been observed within embryonic mesencephalic neuronal grafts more than 10 years after the grafts were therapeutically transplanted into the striatum of PD patients. Specifically, grafted neurons contained a number of Lewy body-like inclusions that stained positively for alpha-synuclein, indicating that host-to-graft transmission of alpha-synuclein pathology had occurred (Li et al. Nat Med (2008), 14(5): 501-3; and Kordower et al. Nat Med (2008), 14(5): 504-6, both of which are incorporated herein by reference).\nFurther, preformed recombinant alpha-synuclein fibrils and alpha-synuclein oligomers can be internalised by cultured cells and neurons, and the direct transfer of alpha-synuclein from donor to recipient cells with the formation of alpha-synuclein inclusions similar to Lewy pathology has been demonstrated (Danzer et al. J Neurosci (2007), 27(34): 9220-32; Volpicelli-Daley et al. Neuron (2011), 72(1): 57-71; and Luk et al. Proc Natl Acad Sci USA (2009), 106(47): 20051-6; each of which are incorporated herein by reference). Injection of preformed synthetic alpha-synuclein fibrils or Lewy body-like alpha-synuclein containing material extracted from the brains of aged alpha-synuclein transgenic mice into the brains of asymptomatic recipient mice promotes the formation of Lewy body-like pathology in host neurons of the recipient animals along with neurodegeneration and neurological deficits (Luk el al. J Exp Med (2012), 209(5); 975-86; and Luk et al. Science (2012), 338(6109): 949-53; both of which are incorporated herein by reference). Alpha-synuclein containing Lewy body extracts isolated from PD brains inoculated into the substantia nigra or striatum of macaque monkeys and mice is rapidly taken up by host cells (within 24 hours) followed by a slower loss of striatal dopaminergic terminals, with cell loss evident after more than a year (Recasens et al. Ann Neurol (2014), 75(3): 351-62; incorporated herein by reference). Similarly, inoculation of mice with brain homogenates derived from patients with the synucleinopathies DLB or MSA triggers alpha-synuclein Lewy-like pathology in the host mice (Watts et al. Proc Natl Acad Sci USA (2013), 110(48): 19555-60; and Masuda-Suzukake el al. Brain (2013), 136(Pt 4): 1128-38; both of which are incorporated herein by reference). Finally, transfer and transmission of both monomeric and oligomeric alpha-synuclein from the olfactory bulb to interconnected brain structures has been demonstrated in mice (Rey et al. Acta Neuropathol (2013), 126(4): 555-73; incorporated herein by reference).\nPassive immunotherapy approaches with antibodies targeting alpha-synuclein have been tested in numerous preclinical alpha-synucleinopathy mouse models (Lawand et al. Expert Opin Ther Targets (2015): 1-10; incorporated herein by reference). Specifically, a study using a monoclonal antibody directed against alpha-synuclein (9E4) has shown in vivo clearance of alpha-synuclein aggregates and pathology, behavioural motor improvements, and neuroprotective effects (WO 2014/058924; which is incorporated herein by reference).\nFurther studies using passive immunisation of alpha-synuclein transgenic mice developed as experimental models of PD/DLB, with the 9E4 monoclonal antibody have shown the antibody to clear alpha-synuclein pathology, decrease synaptic and axonal deficits, abrogate loss of striatal tyrosine hydroxylase fibres, and significantly reduce memory deficits and motor function impairments (Games et al. J Neurosci (2014), 34(28): 9441-54; Bae et al. J Neurosci (2012), 32(39): 13454-69; and Masliah et al. PLoS One (2011), 6(4): e19338; each of which are incorporated herein by reference). Further it has been demonstrated that passive administration of anti-alpha-synuclein monoclonal antibodies in wild-type mice that were injected intrastriatally with synthetic alpha-synuclein preformed fibrils (pffs) led to robust reduction in Lewy pathology, prevention of dopamine neuron loss in the substantia nigra, and a significant improvement in motor impairments that are manifest in the mouse model after pffs treatment (Tran et al. Cell Rep (2014), 7(6): 2054-65; incorporated herein by reference).\nAdditionally, one of the major challenges associated with treating disorders of the CNS with large molecule therapeutics, such as antibodies, is getting these drugs into the affected tissue. The passage of large molecules into the brain and spinal cord is largely restricted by the blood-brain barrier (BBB). The BBB protects and regulates the homeostasis of the brain and prevents the free passage of molecules into most parts of the brain, thereby limiting the treatment of many brain diseases. Transport of essential molecules such as nutrients, growth factors, and hormones is achieved via a series of specific transporters and receptors that regulate passage across the brain endothelial cells. The delivery of biologics and other drugs to the brain therefore represents a significant challenge. Additionally, transport mechanisms appear to exist that rapidly remove antibodies from the brain, presumably to prevent inflammatory responses due to engagement of Fc with effector ligands that promote a pro-inflammatory response.\nOver the last decade, reports of antibody transport across the BBB have emerged where binding to the extracellular domain of the transporter molecules facilitates transcytosis of the receptor antibody complex across the endothelial cell layer.\nThe BBB is mainly comprised of brain capillary endothelial cells, which have specialized characteristics, such as tight junctions, to limit transport of molecules into the brain (Reese et al. 1967, J. Cell Biol. 34: 207-217; Brightman et al. 1969, J. Cell Biol. 40: 648-677; Rubin et al. 1999, Ann. Rev. Neurosci. 22: 11-28), although other cell types, such as pericytes, astrocytes, and neuronal cells, also play an important role in the function of the BBB. Typically, less than 0.1% of a peripherally dosed antibody reaches the brain (Boado et al. 2010, Mol. Pharm. 7: 237-244, Pepinsky et al. 2011, Nat. Neurosci. 8: 745-751). The BBB functions as a physical, metabolic and immunological barrier (Gaillard et al. 2003, Microvasc. Res. 65: 24-31).\nAntibody transport across the BBB can be enhanced by triggering receptor mediated transcytosis on brain endothelial cells. Through this process, engagement of antigens on the luminal side of the endothelial cell can induce the internalization and shuttling of the antibody across the cell, and then its subsequent release into the tissue.\nCurrent drug therapies for PD are mainly focused on treating the motor-related symptoms of the disease. There are currently no marketed or available therapies that can treat or prevent alpha-synucleinopathy.\nAccordingly, there is a need in the art for a therapy for treating alpha-synucleinopathies, particularly in humans."}
{"text": "In cellular networks, e.g., as specified by 3GPP (3rd Generation Partnership Project), increasing traffic demand results in a need for more radio spectrum bandwidth. One way to provide more radio spectrum bandwidth is expansion into unlicensed frequency spectra, e.g., as typically used by WLAN (Wireless Local Area Network) technologies. For example, in 3GPP meeting contribution RP-140240, 3GPP TSG RAN Meeting #63, Fukuoka, Japan, 3-6 Mar. 2014, it is proposed to study extension of the LTE (Long Term Evolution) radio technology for operation in unlicensed frequency bands.\nIn an unlicensed frequency band, typically more bandwidth than the maximum standardized LTE carrier bandwidth (currently 20 MHz) is available. Accordingly, the conventional practice of running all LTE base stations of a network on the same frequency may be suboptimal since the larger available bandwidth allows for reducing intra-cell interference by distributing base stations over multiple different frequency channels. Further, channel quality in unlicensed frequency bands may vary depending on time, location and/or frequency, which means that also the optimum selection of the frequency channel may vary for each base station. Therefore, it is generally desirable to perform selection of the frequency channels for the base stations on the basis of a constantly running automated algorithm.\nIn a general context, the problem of frequency channel selection has been studied intensively for a long time and many different algorithms have been proposed. On a high-level, one may distinguish algorithms that are intended for offline frequency planning of cellular networks like GSM, and real-time algorithms that are intended to be implemented as a Self-optimizing Network (SON) feature in base stations. The main difference between these two categories is the amount if input information they require and the computation time until they deliver results. Offline algorithms are typically allowed to run for a very long time (hours, days) and can afford a significantly higher computational complexity, while real-time algorithms should deliver results in seconds or faster, and may need to cope with limited input information.\nFrequency selection algorithms can be implemented in a distributed or a centralized way. Distributed means that independent algorithm instances run, for example, in each base station. The different algorithm instances influence each other for example in terms of how much interference another instance sees on a given channel. In a centralized approach, all information is gathered in a central location, which allows for a more complete assessment of the overall situation and facilitates finding an optimal solution. Distributed algorithms have a higher risk of being trapped in local minima and are typically iterative, i.e., which means that the system typically runs through a number of suboptimal stages before it may reach a steady state. In a centralized algorithm, even if it is based on iterations, system operation can converge in one step. Further, a centralized algorithm is more likely to find a global optimum because information from various parts of the network can be considered. In each case, finding an algorithm which offers a suitable tradeoff between system performance gains, computational complexity, execution time and other aspects is a complex task.\nA distributed algorithm for WLAN technologies is for example described in “A Self-Managed Distributed Channel Selection Algorithm for WLANs” by D. J. Leith et al., Modeling and Optimization in Mobile, Ad Hoc and Wireless Networks, 4th International Symposium on Modeling and Optimization in Mobile, Ad Hoc and Wireless Networks, pp. 1-9, 3-6 Apr. 2006. This algorithm uses channel selection probabilities based on interference measurements. If measured interference exceeds a threshold, the channel selection probability of this channel is reduced by the same amount. However this algorithm may in some case provide insufficient efficiency in terms of finding the best channel and speed of convergence.\nAccordingly, there is a need for techniques which allow for efficiently controlling frequency channel utilization in a cellular network."}
{"text": "In certain types of building construction, e.g., retrofits, cinder block, or the like, it may not be possible or practicable to run electrical cabling (such as high voltage power lines and low voltage data lines) through the building's walls. In such cases, modular raceway or conduit assemblies are often used to house and route cabling along a wall or other surface. A typical raceway assembly includes a linear or elongate housing having at least one interior passageway that accommodates a length of electrical cabling. (Such a passageway is referred to herein as a “wireway.”) The housing is attached to a wall, and then the cabling is disposed in the interior of the housing. To cover a given span, multiple segments of housing are deployed in an end-on-end manner.\nIn a typical raceway or conduit installation, sections of conduit must be cut precisely to match the wall span to be crossed by the raceway. Such cutting operations, e.g., through sheet metal, are time consuming for installers. Additionally, because the sections are custom cut, it is difficult to adjust or compensate for variations in spacing that stem from adding junction boxes, retrofit outlets, and the like."}
{"text": "Conveyor systems utilizing rollers, also known as idlers, are widely used in a variety of industrial and commercial applications to transport materials. In these environments rollers can be subjected to harsh conditions and considerable use which can cause the rollers to deteriorate rapidly and/or fail prematurely. This is particularly significant since roller failure can cause the shutdown of the conveyor as well as damage to the other components of the conveyor.\nConveyor rollers currently known in the art have problems with such things as establishing an effective seal against contaminants entering the bearings, bearing end play and preload adjustment, shaft and bearing alignment and bearing lubrication. They can also be awkward to handle, very expensive to manufacture and maintenance can often be difficult. There is, therefore, a need for improved conveyor rollers especially for use in industrial applications.\nThis background information is provided for the purpose of making known information believed by the applicant to be of possible relevance to the present invention. No admission is necessarily intended, nor should be construed, that any of the preceding information constitutes prior art against the present invention."}
{"text": "To improve the mechanical and physical properties of castings, it is well known in the art to use inserts which can locally improve the strength, wear resistance or other characteristics of the casting. The inserts are placed within the mold prior to pouring the molten casting material and, upon cooling of the molten casting material, the inserts form an integral part of the finished cast product. A common example of such an application is in the engine block of an internal combustion engine. In particular, aluminum alloy engine blocks, typically an aluminum-silicon alloy, often make use of cast iron inserts, or liners, which are more durable than the cast aluminum walls of the engine block, particularly when aluminum pistons are used. A process well known to those skilled in the art is referred to as the \"Al-fin\" process, which entails dipping a cast iron insert in molten aluminum prior to placing the insert in a mold and casting a molten aluminum around the insert.\nHowever, a disadvantage with the use of cast iron liners is the significant additional weight which is incurred. In addition, cast iron does not conduct heat away from the cylinder as well as aluminum, which raises the temperature of the cylinder and imposes higher temperature-related stresses and wear on the engine's internal components. Another disadvantage with using iron is that there is a mismatch between coefficients of thermal expansion between iron and aluminum and its alloys, which can cause debonding of the insert.\nAs a result, it is generally preferable to provide inserts which are lower in weight while also providing better heat transfer capability and a more closely matched coefficient of thermal expansion. Naturally, aluminum-base alloys and composites are generally suitable in terms of weight, heat transfer and thermal expansion, for use with aluminum castings. Unfortunately, aluminum-base inserts do not metallurgically bond well to casting materials because the insert forms an aluminum oxide layer at the insert's surface. The presence of oxides produces a weak bond because of the inability of the molten casting material to wet the insert's surface.\nTo overcome this problem, one approach known in the art has been to form an insert which can be penetrated by the casting material under high pressure to form a mechanical/metallurgical bond between the casting material and the insert. One such approach uses alumina and carbon fibers which have been highly compressed to form a cylindrical insert. An aluminum casting material is then pressurized sufficiently during the casting process to penetrate the fiber inserts without structurally damaging them. While durability is improved, manufacturing costs are significantly higher than that of iron liners in aluminum blocks.\nAn approach for promoting a metallurgical bond between the insert and the casting material is taught by U.S. Pat. No. 4,687,043 to Weiss et al. Weiss et al. provide an aluminum composite insert whose outer surface is covered with an aluminum alloy. The insert is then coated with a molten solder alloy. The molten solder alloy is selected to have a melting temperature which is below the melting temperature of the insert's aluminum alloy cover layer. The insert is dipped into the molten solder alloy to separate from the cover layer the oxides already present and to prevent the formation of new oxides thereon. In addition, Weiss et al. teach that the casting material must be at a temperature which is higher than the melting temperatures of the insert's cover layer and the solder alloy. This enables the casting material to flush the molten solder alloy from the cover layer during the casting process to expose the cover layer to the casting material, allowing the casting material and the cover layer to form an oxide-free metallurgical bond. The molten solder alloy is intended to be mixed with the casting material and not remain on the surface of the insert.\nIn addition, Weiss et al. teach a zinc solder alloy which contains about 10 to 30 weight percent tin and about 5 to 25 weight percent cadmium to reduce the melting temperature of the solder alloy below the temperature of the casting mold. A disadvantage to the use of tin for the purposes taught by Weiss et al. is that tin embrittles the solder alloy and the interface between the insert and the casting, while also reducing their corrosion resistance. The presence of tin also slows down the age hardening by reducing the Guinier-Preston (GP) zones formation rate, potentially causing a weak interface between the insert and the casting. Similar to tin, cadmium reduces corrosion resistance. But most importantly, cadmium poses an environmental concern in that it is highly toxic. As a result, the use of cadmium is always avoided where possible, and sometimes prohibited.\nThus, it would be desirable to provide a method of promoting a metallurgical bond between a strong, wear-resistant insert and an aluminum alloy casting material which would be economical for use in mass production, while also avoiding the concerns for embrittlement and toxicity."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a vehicle heat exchanger arrangement. More specifically, the present invention relates to recovering power from a cooling fan of a vehicle heat exchanger arrangement.\n2. Background Information\nIn recent years, hybrid vehicles have become more popular. A hybrid vehicle includes an internal combustion engine and an electric motor or a battery as power sources for providing a driving force to a vehicle drive train. With hybrid vehicles, there is an increased need for generating electrical power. When the internal combustion engine is used, the vehicle will generate considerable waste heat during operation. This heat from engine needs to be dissipated or passed to the environment. The engine's coolant is one principle carrier of heat to the environment via a radiator typically placed forward in vehicle. In order to augment heat exchange, particularly at low vehicles speeds, a cooling fan is provided to increase airflow through the radiator. The cooling fan(s) has been historically located at the exit of the heat exchanger (behind the radiator). In some limited applications, a secondary fan is applied upstream of the heat exchanger to supplement airflow and heat transfer. Typically in a vehicle, this cooling action consumes electrical power of about 300 W during operation.\nIn view of the above, it will be apparent to those skilled in the art from this disclosure that there exists a need to recovery the power needed to operate a cooling fan to cool a heat exchanger, such as a radiator. This invention addresses this need in the art as well as other needs, which will become apparent to those skilled in the art from this disclosure."}
{"text": "A large percentage of the oil in petroliferous formations is held within the rock of the formation by the surface forces between the rock, the oil therein, and the aqueous liquid in the formation. As a result, a substantial portion of such oil usually remains in the rock of the formation even when wells traversing the formation are no longer productive. Various secondary techniques, such as thermal recovery, gas injection, and waterflooding have been suggested for the recovery of this fixed oil which remains in the formation after primary recovery methods are no longer feasible to produce additional oil. As a secondary recovery technique, waterflooding is quite commonly employed, and a multiplicity of methods have been suggested for improving the efficiency and economy of oil recovery by the practice of waterflooding. Such methods frequently include incorporation of a water-soluble surfactant in the waterflood liquid, which is injected into the petroliferous formation.\nMore recently, tertiary recovery methods have been used to recover any oil still remaining after the conclusion of the secondary recovery. There are many such methods. A typical method uses an aqueous surfactant followed by the addition of a polymer solution (such polymers are well-known in the art).\nMany types of materials have been suggested as the surfactant for both secondary and tertiary recovery methods. For example many types of sulfonates have been proposed. U.S. Pat. No. 3,302,711 teaches sulfonates derived from a hydrocarbon residue containing 1 to 50 carbon atoms.\nI have found that sulfonates prepared from linear alkanes containing from 20 to 24 carbon atoms provide improved results. The results are better than those obtained on a C.sub.18 alkane sulfonate and much better than that obtained on a C.sub.16 alkane sulfonate.\nHeretofore it has been thought that one property required in a sulfonate for use in surfactant waterflooding was a low interfacial tension between an aqueous surfactant-containing phase and the particular oil. Typically, workers in this field have striven for interfacial tension (IFI) values of 10.sup.-3 dynes/cm. or less. I have found that a surfactant having an IFI of 0.01 to 0.05 dyne/cm. provides excellent results. This is surprising.\nMy invention is suitable for use in both secondary and tertiary waterflood method which uses a surfactant."}
{"text": "As an example of processes of manufacturing a semiconductor device, a process of forming a film such as a nitride film, an oxide film or the like on a substrate is often carried out by supplying, e.g., a precursor containing silicon and a reactant such as a nitriding gas, an oxidizing gas or the like to the substrate.\nIn some cases, depending on the kinds of precursors, when performing the aforementioned film forming process, it may be difficult to improve a deposition rate, and then the productivity of a film forming process deteriorates."}
{"text": "Inkjet printers find uses in a wide range of applications. Reductions in ink drop application sizes have made inkjet printers useful in color printing, such as the printing of photographs.\nWhen rows or columns of nozzles are used to eject drops to form an image, small horizontal and/or vertical bands may be created by a plugged or malfunctioning nozzle. Horizontal and/or vertical bands also may be caused by directionality errors in ejected drops. Depending on manufacturing variations in the printhead, ejected drops may not always be ejected exactly perpendicular to the print medium. The bands created by such plugged nozzles, malfunctioning nozzles, and/or ejection directionality may be detected by the human eye thereby diminishing the quality of the printed image.\nThere are methods that may be used to detect nozzles that are not working properly. Such methods may be fairly expensive to implement in a consumer product. If a non-functioning nozzle is detected, compensation may be made by passing another working nozzle over the portion of the image associated with the non-functioning nozzle. However, high-speed printing may be done with only one or two passes of the nozzles over the same location on an image that is the printed. This makes it difficult to compensate for non-functioning nozzles. This is especially difficult for page wide array printheads. Since the printhead of a page wide array does not move, there may be no opportunity to use another nozzle to compensate for the plugged or malfunctioning nozzle."}
{"text": "Conventionally, as an AC generator for a vehicle, there is known the one which employs a Lundell-type rotor including a pair of claw-like magnetic poles axially supported by an intermediate portion of a shaft (for example, see Patent Literature 1).\nIn the rotor, a slip ring device including a molded-body main body formed by molding, which retains a first ring and a second ring disposed coaxially so as to be separated away from each other in an insulated fashion, is press-fitted over a distal end portion of the shaft.\nOn the other hand, in the case where a molded body including the molded-body main body is formed by molding, a weld line is generated when branch flows of a melted resin meet inside a die. In order to minimize the generation of the weld line, measures to increase temperatures of the melted resin and the die and to increase an injection pressure are taken (for example, see Patent Literature 2)."}
{"text": "Automatically identifying the locations of objects and their parts in video is important for many tasks. For example, in the case of human body parts, automatically identifying the locations of human body parts is important for tasks such as automated action recognition, human pose estimation, etc. Body parsing is a term used to describe the computerized localization of individual body parts in video. Current methods for body parsing in video estimate only part locations such as head, legs, arms, etc. See e.g., “Strike a Pose: Tracking People by Finding Stylized Poses,” Ramanan et al., Computer Vision and Pattern Recognition (CVPR), San Diego, Calif., June 2005; and “Pictorial Structures for Object Recognition,” Felzenszwalb et al., International Journal of Computer Vision (IJCV), January 2005.\nMost previous methods in fact only perform syntactic object parsing, i.e., they only estimate the localization of object parts (e.g., arms, legs, face, etc.) without efficiently estimating semantic attributes associated with the object parts.\nIn view of the foregoing, there is a need for a method and system for effectively identifying semantic attributes of objects from images."}
{"text": "The present invention relates to a silver halide photographic material. More particularly, the present invention relates to a silver halide photographic material capable of providing a high-contrast and sensitivity photographic image in a rapid and consistent way.\nHigh-contrast photographic image is usually employed in forming characters or halftone image in photomechanical processes, or forming fine-line image in ultra-high-precision photomechanical processes. Certain kinds of silver halide photographic materials are known to be capable of forming photographic image with very high contrast. In one conventional method, a light-sensitive material composed of a silver chlorobromide emulsion that has an average grain size of 0.2 .mu.m, that has a narrow grain size distribution, that comprises uniformly shaped grains and that has a high silver chloride content (at least 50 mol %) is processed with an alkaline hydroquinone developer of low sulfite ion concentration to form an image of high sharpness and high resolution, for example, a halftone image or a fine-line image. Silver halide light-sensitive materials of a type to be processed by this method are known as \"lithographic\" light-sensitive materials.\nThe photomechanical process involves a step of converting the original of continuous tone to be reproduced into a halftone image, or a step in which the continuous change in the density of the original is converted to a set of dots having areas proportional to the varying densities. To form a halftone image, the original is imaged onto a \"lithographic\" light-sensitive material via a crossline screen or a contact screen and subsequently developed. However, development of a lith-type light-sensitive material (i.e., a silver halide photographic material containing a silver halide emulsion consisting of fine grains of uniform size and shape) with an ordinary black-and-white developer results in the formation of a lower quality image than when it is developed with a lith-type developer. Thus, the lith-type light-sensitive material is conventionally processed with a lith-type developer that has a very low concentration of sulfite ions and which contains hydroquinone as the sole developing agent. However, developers of the lith type are very low in storage stability because they are highly sensitive to autooxidation. Hence, there is a strong need for controlling them to provide consistent quality of development in their continuous use and substantial efforts have been made to improve the storage stability of lith-type developers.\nOne of the methods proposed so far to attain object is replenishment with two separate liquids, one being a replenisher intended to compensate for the deterioration in developer's activity during development (compensation for process fatigue) and the other being a replenisher intended to compensate for the oxidative deterioration over time (compensation for aging fatigue). This method of using two types of replenishers is commonly adopted with automatic processors in photomechanical processes. However, because of the need to control the balance in replenishment with two liquids, this method suffers the problem of complexity in both equipment and operations.\nAnother problem with the use of lith-type developers is that they have a long induction period (i.e., the time required for an image to appear upon development) and hence are incapable of providing rapid access to image.\nMethods are known that are capable of rapid formation of high-contrast image without using developers of the lith type. They are characterized by containing hydrazine compounds in silver halide light-sensitive materials [see U.S. Pat. No. 2,419,975 and JP-A-51-16623 (the term \"JP-A\" as used herein means an \"unexamined published Japanese patent application\"), JP-A-51-20921, etc.]. According to these methods, the concentration of sulfite ions in developers can be held high enough to accomplish development with their storage stability being maintained at enhanced levels. However, in order to obtain a contrasty image by these methods, a developer with fairly high pH is necessary, which presents a problem in keeping the developer stable to produce a photographic image rapidly and with high sensitivity. The high pH of the developer tends to cause fogging and to prevent it, various organic restrainers must be contained in high concentrations but then sensitivity is compromised. A further problem is that a contrasty image is susceptible to defects such as pinholes if dust particles or other foreign substances are deposited on the light-sensitive material."}
{"text": "The present invention relates generally to manual control devices for automotive vehicles.\nApplicant is aware of a number of both patented and unpatented devices for manually controlling the brake and/or accelerator pedals of automotive vehicles. The inventors and the respective patent numbers of which applicant is aware of as follows:\nPawl 3,275,093\nBhattacharya 4,143,734\nAppley 4,228,865\nDowden, et al. 4,438,835\nJohnson, et al. 4,476,954\nAhnafield 4,722,416\nA number of unpatented devices are also known to applicant. One such device is distributed by the Ricon Corporation of Sun Valley, Calif., which activates the throttle of the accelerator and brake pedals by a push-pull motion of the longitudinal handle. A very similar device is marketed by Manufacturing and Production Services Corporation of San Diego, Calif. Applicant is also aware of the hand control device manufactured and marketed by Handicaps, Inc., which activates responsive rods to forward pivotal movement and downward pivotal movement of the elongated control arm, respectively.\nApplicant's invention is directed to a unique vehicle hand control device which activates the accelerator pedal by rotation of a control arm about its longitudinal axis via manual corresponding rotation of a uniquely configured handle. Separately, the brake pedal of the vehicle is activated by a forward pushing motion of the handle which pivots the control arm about its pivotal connection to a supporting control arm tube positioned adjacent the vehicle's steering column."}
{"text": "About one third of one's lifetime is spent sleeping and a not insignificant effort has been spent attempting to achieve comfort during this time period. Many inventors have attempted to approximate the comfort that one experiences in floating in a body of water such as the Great Salt Lake in Utah. It is believed that the California Indians spent the better part of cold winter days semi floating in pools of hot mud fed by the natural warm springs. The so called \"water bed\" has become popular, especially in the State of California. Water beds have taken two approaches; those that provide a \"floatation\" or \"buoyancy\" effect, and those that result in a \"hammock\" effect. The first efforts known to Applicant to provide a \"floatation\" effect is disclosed in White, U.S. Pat. No. 184,487, Nov. 21, 1876. White disclosed a plurality of water tight sacks instead of a \"single rectangular oblong air and water tight sack\". The White bed consisted of a rigid perimeter and the entire bed was filled with the water sacks. The rigid border probably provided a \"floatation\" effect. On June 15, 1971, Hall received U.S. Pat. No. 3,585,356 for a water bed consisting of a single flexible substantially inelastic bladder contained by a perimeter rigid frame. Hall, appears to be the first to fully explain the \"floatation\" effect which results from the use of a frame which provides lateral support to the flexible inelastic bladder. Hall filled the entire bed with water resulting in an extremely heavy piece of furniture. The entire body of the person was supported on the inelastic water bladder.\nThe \"hammock\" effect water bed is discussed in the Labianco patent U.S. Pat. No. 3,840,921. Labianco taught that if a person lies on a pillow type water bag which is not supported on its sides, the top of the water bag is placed in tension and \"it will create conventional bed pressure points, and will not conform and adjust properly to the different weight proportions of the upper and lower torso of the user.\" In effect, the body is supported in a \"hammock\" type attitude.\nBoth the laterally supported \"floatation\" beds and the \"hammock\" type non-laterally supported beds have the problem of excessive weight. More recently a great deal of activity has centered on the problem of reducing weight by providing air or foam rubber or plastic perimeters. Examples of such reduced weight beds are taught by Tobinick in his U.S. Pat. No. 3,702,484, U.S. Pat. No. 3,789,442 and 3,815,165; and Tinnel U.S. Pat. No. 4,015,299. All of these patents lack sufficient structure to provide the necessary lateral support to the flexible water container to result in a \"floatation\" type support. All of these beds which are combinations of unconnected foam materials, water containers and support sheathes provide a \"hammock\" type of support.\nReduction of mattress weight by replacing portions of the water area with lightweight foam materials has created a comfort problem in that a person lying on beds constructed in the manner of Tobinick, U.S. Pat. No. 3,789,442 and Tinnel, U.S. Pat.No. 4,015,299 can feel the edge of the foam material where it borders the cavity holding the flexible water bladder. This problem can be overcome, but not with foam mattresses which have cavities which do not provide sufficient lateral support. Hall, supra, taught that lateral support about the perimeter of the water bladder was essential. Without lateral support, a body resting therein sinks to an undesirable depth in the water. If the legs are resting on the foam at an elevation much higher than the heavier portions of the body such as the hips or shoulders, it gives the person an uncomfortable feeling and in addition he will feel the boundry between the foam and the cavity holding the water bag. This disparity in vertical displacement is even accentuated in Tobinick, U.S. Pat. No. 3,789,442 because his cavity walls are slanted.\nFinally, the prior art mattresses having foam material areas and water supported areas failed to recognize the great difference in evenness of body weight and the pressure points which occur at the hips and shoulders. Under these pressure points, foam material is more compressed than it is under the arms or legs. This unevenness in support can cause discomfort over protracted periods of time unless the sleeper shifts his weight. This unevenness of support is especially apparent in Tobinick U.S. Pat. No. 3,789,442 in which the hips and legs are supported on the water bag and the shoulders and arms are supported on foam material. While the Tobinick hospital bed had the advantage that the upper portion can be raised from the horizontal without disturbing the water cavity, the level of comfort in the hips and the shoulders is quite different."}
{"text": "A trailer, such as a travel trailer, boat trailer, utility trailer, etc., usually has a tongue portion at the front of the trailer with a coupling that connects the trailer to a hitch on the towing vehicle. Safety chains are typically attached permanently to the tongue at one end of each chain and have hooks at the other ends of the chains to removably connect the chains to the hitch. An electrical wiring harness is typically hard-wired to the trailer tongue at one end and includes a plug at the other end to removably connect the harness to an electrical connector on the towing vehicle. These chains and harnesses present a storage and clutter problem when the trailer is not in use.\nOne proposed solution to the storage problem is disclosed in U.S. Patent Application Publication No. 2006/0061064 which discloses a protective cover for a trailer tongue. The cover is essentially a bag that fits over the coupling and part of the tongue and allows the user to insert the safety chains and wiring harness in the bag. Placing everything in a bag, however, is not a satisfactory solution. The chains and/or harness can become tangled in the bag, and moreover, it does not provide a sense or organization.\nAnother aspect of trailer storage relates to tongue jacks. Most trailers have a tongue jack with an extendable column that raises or lowers the trailer tongue to facilitate coupling the trailer to the towing vehicle, and also enabling the trailer to be leveled when it is parked. Some tongue jacks have manually operated cranks to raise and lower the tongue, while others have electric or pneumatic motors. An electric tongue jack typically has a head with one or more motors, actuators, buttons, switches, displays, lights, etc. at the top of the extendable column.\nIn recent years, electric tongue jacks have become more sophisticated, functional, and stylized. An example of such a jack is shown in FIG. 1 which illustrates a Smart Jack by Lippert Components, Inc., (LCI). The jack 10 of FIG. 1 includes a power head 12 that contains a motorized lift mechanism, a telescoping support column 14, a flange 16 to attach the jack to the frame of a trailer, and an adjustable foot 18. The power head 12 has an electronic control panel 20 that includes buttons and displays for controlling the operation of the jack.\nIn addition to functionality, these jacks provide their owners with a sense of enjoyment and pride in the fine design and fabrication that goes into these devices. They can be very expensive, and owners may therefore want to preserve and display their investment."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an oral hygienic device and more specifically a disposable toothbrush containing medicated or non-medicated dentifrice and either a single dose or multiple doses of dentifrice to enable a single use and disposal or disposal after multiple uses. The disposable toothbrush includes a tubular housing of generally oval-shaped configuration and a slidably extendable and retractable, longitudinally elongated bristle head carrier having a pivotal bristle head at one end which partially telescopes into the end of the tubular housing when the bristle head and bristle head carrier are in retracted position and pivots to an angulated relation to the bristle head carrier when in the extended position. The oval-shaped housing is approximately 3\" in length with the carrier and brush head extending approximately 3\" from the housing when in extended position to enable the housing to form a handle for the toothbrush. The oval-shaped transverse configuration of the housing includes a major diameter of just over 1\" to enable the device to be easily carried in a pocket when in a retracted or collapsed position without producing a noticeable bulge or distortion of a garment. A latch structure retains the brush head carrier in extended position and the brush head telescopes into the end of the tubular housing when in retracted position. In the single use model of the toothbrush, the dentifrice is preloaded in the bristle head and in the multi-use model, there is a reservoir for a multi-load dentifrice cartridge in the tubular housing which forms the handle to enable the user to reload the bristle head with dentifrice for up to five use cycles. The bristle head includes a snap-on feature to enable the bristle heads to be replaced during each use cycle.\n2.Description of the Prior Art\nIt is well known that pathogenic germs thrive on the bristles of a conventional toothbrush and it has been established that the toothbrush is responsible for the transmission of various infections. Various pathogenic organisms such as bacteria, fungi, yeasts and viruses can remain alive on a moist toothbrush for up to a week. The toothbrush is used in the mouth which is a known bacteria containing area and the usually moist bathroom in which the toothbrush is stored assists and in some instances accelerates microbial growth. Cleaning the conventional toothbrush by a simple \"rinse and tap method\" is generally ineffective in removing toothpaste residue, saliva, and food debris trapped in the bristles or tufts all of which provide an environment to maintain pathogenic organisms alive which can provide a potential for cross-contamination caused by the various pathogenic organisms remaining alive on the toothbrush.\nIt is also recognized that it is desirable to brush your teeth at more frequent time intervals, such as after eating. However, this is difficult to accomplish since individuals usually do not have toothbrushing equipment available throughout a work day, when travelling or the like. There have been efforts made to provide a collapsible or travelling toothbrush which are sometimes provided with a supply of dentifrice. However, such devices are relatively bulky and not disposable. Thus, while such devices facilitate brushing the teeth at more frequent time intervals, they still have the problem of pathogenic organisms remaining alive on the toothbrush bristles.\nThe following U.S. patents are relevant to this invention. U.S. Pat. Nos.\n2,454,995, 2,620,500, 2,668,973, 2,828,504, 3,205,521, 4,344,535, 4,482,263, 4,530,129, 4,542,828, 4,695,177, 4,866,809,\nThe prior art in this field of endeavor does not include a disposable toothbrush of small compact size which can be extended to approximate a full size toothbrush with the bristle head being angled with the toothbrush being constructed in a manner to render the device economically disposable after a single use or after a limited number of uses."}
{"text": "1. Field of the Invention\nThis invention pertains to the field of bit-mapped alphanumeric and graphic processors, or bit-mapped raster scan video controllers, and in particular to the logic and circuit necessary to implement a variety of strategies for sharing access to display memory among a number of processes under the programmable control of the raster scan video controller. The invention is useful for high and low performance CRT systems, black/white or color, especially those capable of accessing display memory as needed to create and update an image on a video display.\n2. Description of the Prior Art\nMost presently available video display systems typically include a processor, a video controller, a display memory containing a single current screen image, other system memory, and a raster scan video display. In normal (steady-state) operation, the video controller continually reads out the contents of the display memory and transforms the information read out to the signalling necessary to control the raster scan beam while it is in its active display time. The video controller also provides the horizontal and vertical retrace signalling at appropriate intervals, and blanking of the raster scan beam during retrace.\nThe processor also has access to the display memory, so that it can change the current screen image. This access may be \"through\" the video controller or \"around\" it. The subject invention applies to the former type of system. In either case, use of the display memory is typically carefully controlled between updating and display accesses, to prevent breakup of the video image while it is being changed.\nIn bit mapped display systems, display memory must be accessed by or shared among (1) the display process that keeps the image on the CRT and one or both of: (2) a dedicated hardware engine that updates or changes the image; and/or (3) a microprocessor that updates or changes the image. Existent CRT controllers typically adopt a fixed strategy for the sharing--either alternating accesses between (1) and (2) and/or (3), and/or allowing 2 and/or 3 access during retrace or blanking times. In most prior art systems, the allocation of display memory accesses between the display process and updating is external to the video controller. The two types of accesses are kept physically separate by logical circuits such that the allocation is fixed and not subject to change.\nDepending on the timing of the various parts of the system, the display memory may be available for updating a) only during vertical retrace periods, or b) during horizontal and vertical retrace periods, or c) during retrace periods plus alternating memory cycles during the active display time of scan lines. In any of these cases, however, updating of display memory typically proceeds at a rate slower than could be achieved without interference from the video controller's display accesses.\nThe invention is a mechanism by which a raster scan video controller can be programmed for any of a variety of sharing strategies, in accordance with application requirements.\nAnother object of the present invention is to improve performance of a raster scan display system because the update and video operations can occur simultaneously in many cases."}
{"text": "This invention concerns filters and more particularly filters of the type adapted to remove suspended solids from liquids, such as for dewatering or desludging applications.\nSuch applications usually involve large volumes of liquid to be filtered, and large filtering areas are necessitated. Typically, pressurized flow is necessary to provide adequate flow rates, and a high degree of automation of the filter operations is highly desirable.\nA known arrangement for conducting filtration for such applications is a vertically stacked pressurized chamber filter arrangement. A typical arrangement of this type is shown in U.S. Pat. No. 4,346,003. In this arrangement, a plurality of filter housings are mounted in a vertical array, with each filter housing defining a filter chamber which is pressurized with the liquid to be filtered. The filter flow passes through a filter media taking the form of an endless belt which has a segment passing into each filter chamber, and out around rollers and back into the next filter chamber, and so on, through the vertical array. Each filter consists of separable housing members which are simultaneously moved apart vertically to allow advancement of the belt segments from within the filter housing. This movement is necessary to remove the accumulated filtered solids and clean the belt as the filtration process proceeds.\nThe necessity for separating the filter housings requires a complex and very sturdy mechanism which must withstand the enormous forces involved with large area filtration under pressurized conditions.\nIn related application Ser. No. 405,524, filed Aug. 5, 1982, now U.S. Pat. No. 4,430,231, there is disclosed an arrangement for sealing openings into a filter housing by means of inflatable hose seals, such that the housing members may be permanently closed and only the inflatable sealing apparatus operated to enable the filter belt and accumulated solids to be periodically moved out of the filtering cavities and for disposal.\nThis arrangement eliminates the need for the complex mechanical appratus for opening and closing the housing components.\nHowever, a relatively complex belt drive mechanism is required since each segment of the endless belt is passed into and out of each housing successively into a next succeeding housing in the vertical array. Thus, a very great length of filter belt is involved, and a complex roller and drive arrangement necessitated to advance all of the belt segments simultaneously.\nAn effective sealing of the filter belt edges in the arrangement described in the aforementioned patent application is more difficult since flow around and under the edges of the filter is more likely. Any build-up of sludge will of course necessitate disassembly of the unit for periodic clean out. This would be highly disadvantageous. In the separable housing-type filters, complete sealing of the belt lateral edges may be achieved by sealing engagement of the housing.\nAn alternative approach to the endless belt has involved the use of individual rolls of a disposable flexible sheet filter media, which are individually associated with each housing in the vertical array. Since the media is disposable, it need not be circulated back into a succeeding filter chamber, and may pass therethrough for collection and disposal on the opposite side of the housing chamber. However, the cost of such disposal media may become substantial for typical high volume applications of the sort described.\nAdditionally, the disposable media is typically of much lighter weight and is less easily driven through the housing, particularly when loaded with the accumulated filtered solids. If the media passes atop a perforated plate or grid, filtration pressure forces the disposable media into openings in the support plate, and a high pulling force is necessary. This has led to difficulties in reliably powering the movement of the filtering media out of the filtration chamber, in addition to incurring the cost of the filter media material.\nAnother problem has been encountered in such vertically stacked pressure filter units when filtering liquids containing very fine solids. A significant accumulation of very fine solids passing through the media may occur in the collection chambers associated with the filter unit. Accumulation of these solids in the collection chambers may necessitate dismantling and disassembling of the entire apparatus in order to clean out the accumulated fine solid materials.\nA further problem has been encountered in pressurized chamber filters using a flexible sheet media. That is, the media must be supported on a plate which provides support for the filter media to resist the pressure differentials required for filtration.\nThe solid areas of the support plate reduces the area of the permeable filter media available for filtering flow since the media tends to be sealed against such solid areas. On the other hand, a reduction in the solid areas weakens the support plate, and at relatively higher pressures, the structural integrity of the support plate may be threatened.\nAccordingly, it is an object of the present invention to provide a filter of the type including an endless belt member which passes into a filtration chamber, but which does not require a complex belt drive arrangement.\nIt is a further object of the present invention to provide such filters in which a plurality of units may be stacked in a vertical array to establish independently operated filter belt drive arrangements.\nIt is still a further object of the present invention to provide an edge sealing arrangement for the filter of the type described in which the belt passes into the interior of the filter chamber of a permanently assembled filter housing through selectively sealed openings.\nIt is yet another object of the present invention to provide an arrangement for collecting the filtered liquid flow which allows for clean-out of the collection space, without disassembly of the filter apparatus, particularly in a stacked vertical array of the filter units.\nIt is another object of the present invention to provide a support plate configuration which allows for maximum filtering flow area but which provides sufficient structural integrity to resist the stresses induced by high pressure differential filtration, yet which is relatively low in cost to manufacture."}
{"text": "The disclosed embodiments of the present invention relate to power amplification, and more particularly, to a multi-stage digitally-controlled power amplifier.\nA polar transmitter offers some advantages, such as a potential for reducing complexity and current consumption in the modulator path as well as eliminating the problem of image rejection, thus the polar transmitter is more suitable for implementation in advanced complementary metal oxide semiconductor (CMOS) processing technologies. More specifically, the polar transmitter is a transmitting device that splits a complex baseband signal explicitly represented by an amplitude-modulated (AM) contented component and a phase-modulated (PM) contented component, instead of an explicit in-phase component and a quadrature component. These two orthogonal components are then recombined into a radio-frequency (RF) output signal to be transmitted over the air.\nAn all-digital radio-frequency (RF) transmitter front-end circuit may be employed by the polar transmitter to enhance power efficiency, reduce the hardware cost and reduce the chip size. One conventional implementation of the all-digital RF transmitter front-end circuit is a digitally-controlled power amplifier (DPA), acting as an RF digital-to-analog converter (RF-DAC). The DPA may include a plurality of DPA cells for combining the AM signal and the PM signal and delivering an integral signal having a desired RF carrier frequency and a required power level. Therefore, how to arrange and control these DPA cells to achieve the desired DPA functionality is a significant concern in the pertinent field.\nRegarding a conventional design of a DPA cell, a high-efficiency switching-mode (inverse class-D/class-E) power amplifier is employed. However, current switching-mode DPAs are inductor-loaded. Thus, due to the inherent characteristics of the inductor-loaded switching-mode amplifier, the voltage swing will be greater than π×VDD, where VDD is the supply voltage. To reduce the voltage swing, a lower supply voltage VDD must be used. For example, a DC-DC converter and/or a low-dropout (LDO) regulator may be used to convert a high DC voltage provided by a battery into a low supply voltage VDD needed by the DPA. Unfortunately, the battery efficiency would be degraded due to the power conversion.\nBesides, the AM sampling using a periodic sampling clock would introduce out-of-band (OOB) noise/replica. One conventional solution to reduce the OOB noise/replica is to use a higher sampling rate. However, this would lead to larger power consumption. Another conventional solution is to use a higher digital-to-analog converter (DAC) resolution. However, this would have layout and physical limit. Yet another conventional solution is to use an RF bandpass filter. However, such an RF bandpass filter has a low quality factor and consumes a large area.\nThus, there is also a need for an innovative DPA cell design which has improved efficiency and can effectively reduce the undesired OOB noise/replica.\nMoreover, due to the feedback path established by magnetic coupling and/or direct coupling (e.g., coupling via the printed circuit board (PCB) ground and/or the package ground), the transmitter output may be fed back to a clock source in the transmitter, which may degrade the transmitter performance. Thus, there is also a need for a pulling mitigation mechanism employed to improve the transmitter performance."}
{"text": "The present invention relates to a method of mounting chips, and particularly to a method of mounting a large number of fine chips, for example, fine semiconductor chips having sizes of the order of m on boards.\nFunctional devices formed of semiconductor chips such as LEDs and ICs having fine structures have been mounted on boards by forming the semiconductor chips into shapes such as trapezoidal shapes by etching while providing recesses corresponding thereto on the board side, wherein the semiconductor chips are mounted in the recesses.\nThe above method, however, has the following problem: namely, in the case of using semiconductor chips made from materials not easy to be formed into trapezoidal shapes by etching or of fine dimensions, the mounting of these semiconductor chips becomes difficult, or the accuracy in shape of the chips and/or the positional accuracy in mounting of the chips are lowered, thereby degrading the mounting yield.\nIn mounting semiconductor chips such as light emitting devices used for a display on a board, the semiconductor chips have been required to be mounted on a board surface having a specific area with a high efficiency and a high yield, and further, in consideration of the cost. Further, along with the tendency toward miniaturization of semiconductor chips, the handling of the semiconductor chips at the time of mounting have come to be difficult. Accordingly, it has been expected to develop the semiconductor chip structures capable of effectively coping with the above terms, that is, the efficiency, yield, cost, and handling at the time of mounting, and to develop the method of desirably mounting these semiconductor chip structures on boards.\nOne method of mounting fine chips such as fine semiconductor chips has been disclosed, for example, in U.S. Pat. Nos. 5,545,291, 5,824,186, and 5,904,545, Japanese Patent Laid-open (Translation of PCT International Publication) No. Hei 9-506742, and Japanese Patent Laid-open No. Hei 9-120943. In this mounting method, each chip is tapered to be determined in its vertical direction, and buried in the corresponding recess formed on a board. In this case, the chips are mixed in a solution such as water or alcohol into a slurry, and the slurry is made to flow on the board.\nAnother method of mounting fine chips, such as IC chips, using static electricity for mounting the IC chips on a board has been described in U.S. patent application Ser. No. 07/902,986. In this mounting method, the chips are vibrated by static electricity and are arranged in a state in which the positional energy of the chips is minimized.\nAccording to the former mounting method characterized by tapering semiconductor chips, although a sufficiently large yield can be kept for semiconductor chips whose sizes are in the range of several hundreds μm or more; however, the yield is lowered for semiconductor chips whose sizes are in the range of several tens μm or less. Accordingly, this method has a problem from the practical utilization.\nThe latter mounting method characterized by using static electricity requires an apparatus for vibrating fine chips with static electricity. Further, according to this method, since fine chips collide with each other by mechanical vibration, there may occur an inconvenience that part of each chip be damaged. From this viewpoint, this method is unsuitable to practical use."}
{"text": "Frequently, users of electronic devices, such as laptop computers, wish to view displayed content at a higher magnification. For example, a user may wish to view a desired region of a displayed document at a higher magnification. In order to allow this, a function whereby a user may “zoom” in to, or view at an increased magnification level, the desired region of the displayed content may be used. However, when viewing the desired region at the increased magnification level, a context of the region within the entire content is difficult to ascertain or may be lost.\nIt is an object of embodiments of the invention to at least mitigate one or more of the problems of the prior art."}
{"text": "The invention described and claimed herein relates to the field of detecting rotation of a structure about an axis. In particular, the claimed invention relates to the field of monolithically-integrated gyroscopes."}
{"text": "The present invention relates generally to integrated circuit memory devices and, more particularly, to a method and system for implementing dynamic refresh protocols for DRAM based cache.\nMemory devices are used in a wide variety of applications, including computer systems. Computer systems and other electronic devices containing a microprocessor or similar device typically include system memory, which is generally implemented using dynamic random access memory (DRAM). A DRAM memory cell generally includes, as basic components, an access transistor (switch) and a capacitor for storing a binary data bit in the form of an electrical charge. Typically, a first voltage is stored on the capacitor to represent a logic HIGH or binary “1” value (e.g., VDD), while a second voltage on the storage capacitor represents a logic LOW or binary “0” value (e.g., ground). A principal advantage of DRAM is that it uses relatively few components to store each bit of data, and is thus a relatively inexpensive means for providing system memory having a relatively high capacity.\nAs a result of the package/pin limitations associated with discrete, self-contained devices such as DRAMs, memory circuit designers have used certain multiplexing techniques in order to access the large number of internal memory array addresses through the narrow, pin-bound interfaces. Because these discrete DRAMs have been in use for some time, a standard interface has understandably emerged over the years for reading and writing to these arrays. More recently, embedded DRAM (eDRAM) macros have been offered, particularly in the area of Application Specific Integrated Circuit (ASIC) technologies. For example, markets in portable and multimedia applications such as cellular phones and personal digital assistants utilize the increased density of embedded memory for higher function and lower power consumption. Unlike their discrete counterparts, the eDRAM devices do not have the limited I/O pin interfaces with associated memory management circuitry. In fact, the typical I/O counts for eDRAM devices can number in the hundreds.\nOne disadvantage of DRAM (including eDRAM), however, is that the DRAM memory cells must be periodically refreshed as the charge on the capacitor eventually leaks away and therefore provisions must be made to “refresh” the capacitor charge. Otherwise, the data bit stored by the memory cell is lost. While an array of memory cells is being refreshed, it cannot be accessed for a read or a write memory access. The need to refresh DRAM memory cells does not present a significant problem in most applications; however, it can prevent the use of DRAM in applications where immediate access to memory cells is required or highly desirable.\nThus, in certain instances, the refresh process involves accessing memory locations corresponding to the same location from which data is needed for system operation. This contention with refresh increases the average latency of the operational accesses. In addition to decreasing memory availability, refresh requirements of eDRAM can also negatively impact power and performance and add complexity to system implementation. On the other hand, alternative refresh solutions (such as conventional selective/partial refreshing, for example) save power, but at the cost of degraded performance from the loss of data. Accordingly, it is desirable to be able to provide improved approaches to refresh protocol management for simultaneously improving memory availability, system performance and power consumption."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a POE (power over Ethernet)-prioritized active splitter, and more particularly pertains to a POE-prioritized active splitter that uses an intelligent POE power supply with a prioritized back-up response, allowing a lower cost UPS (uninterruptible power supply) system to be used, and also providing a longer back-up time period by the UPS system for critical components of a prioritized network having different priorities such as a security system or any prioritized network such as an internet or intranet based network, a monitoring system network, communication system network, etc.\n2. Discussion of the Prior Art\nA recently released standard, IEEE 802-3af for POE (power over Ethernet), addresses the technical field of development wherein electrical power/current is provided to networking hardware devices over Ethernet category 5 or higher data electrical wiring/cables. No additional AC or DC power line is needed to supply electrical power at the product or hardware locations, thereby minimizing the amount of required electrical wiring/cables and/or the problem and expense of installing additional power lines and outlets.\nIn a residential or commercial security system, the security system control panel is frequently connected to various security system peripheral hardware devices using structured wiring/cables. POE devices use category 5 or higher structured wiring/cables, which generally include several twisted wire pairs, for providing both data transmission to and electrical power for the security system peripheral hardware devices.\nA conventional security control panel uses a proprietary bus and separate wiring/cables within a building to communicate with zone-expansion modules, RF (radio frequency) receivers, internet alarm reporting devices, IP (internet protocol-allows a camera to connect to and communicate directly over an Ethernet network) cameras, etc., wherein the separate wiring/cables carry the data communications traffic as well as provide source electrical power to the security system peripheral devices.\nIn most new residential construction and commercial buildings, the building is prewired with category 5 wiring/cables for business communications, or for a home network, typically connected through a router to a cable modem or ADSL (asymmetrical digital subscriber line) connection and then to an internetwork (the internet). While the use of the existing category 5 wiring/cables reduces the time and cost of wiring a security system, it has several associated disadvantages and problems:\nProblem 1) Should the LAN (local area network) become disabled for any reason, then data messages between the various component devices of the security system may be interrupted and lost.\nProblem 2) Should the local AC power supply become interrupted such as by a general power outage, then the security system component devices also lose power, and the security system becomes inoperative.\nProblem 1) is typically addressed and mitigated by providing a second redundant communication system, such that if internet connectivity to a central security station is lost, or the LAN traffic is disrupted, then an independent trouble report message can be sent, perhaps by conventional POTS (plain old telephone service), or by using a long range RF network. This solution requires that the entire security system, including the back-up communication system, remains powered. This is really a poor solution since although a trouble report message may be sent, the security system still remains inoperative until the LAN becomes operative and functional again.\nProblem 2) is typically addressed and mitigated by the use of various power supply back-up measures, for example by providing stand-by generators in the larger systems, or by providing a low cost battery powered UPS (uninterruptible power supply) system in residential systems. However the cost and size of UPS systems, and the required maintenance on these UPS systems has not resulted in widespread use of UPS systems of sufficient capacity, particularly in low cost residential systems."}
{"text": "With the increasing popularity of playing streaming audio and video over networks such as the Internet, there is a need for optimizing the data transferred from a server to a client such that the client's experience is maximized even if network conditions during playback are inconsistent. Optimizing the client's experience involves making encoding decisions such that the video can be transferred and reconstructed with a minimal number of errors.\nThe term “streaming” is typically used to indicate that the data representing the media is provided by a host computer over a network to a playback device (i.e., a media playback computer implemented as any of a variety of conventional computing devices, such as a desktop PC, a notebook or portable computer a cellular telephone or other wireless communication device, a personal digital assistant (PDA), a gaming console, etc.) The client computer typically renders the streaming content as it is received from the host, rather than waiting for the entire file to be delivered.\nThe quality level is generally dictated by the bit rate specified for the encoded audio or video portions of the input stream. A higher bit rate generally indicates that a larger amount of information about the original audio or video is encoded and retained, and therefore a more accurate reproduction of the original input audio or video can be presented during video playback. Conversely, a lower bit rate indicates that less information about the original input audio or video is encoded and retained, and thus a less accurate reproduction of the original audio or video will be presented during video playback.\nGenerally, the bit rate is specified for encoding each of the audio and video based on several factors. The first factor is the network condition between the server and the client. A network connection that can transfer a high amount of data indicates that a higher bit rate can be specified for the input video that is subsequently transferred over the network connection. The second factor is the desired start-up latency. Start-up latency is the delay that a video playback tool experiences when first starting up due to the large amount of data that has to be received, processed, and buffered. Start-up latency can also occur after a seek operation, where the user selects variable positions in the streaming content to view. A third factor is the processing capabilities of the playback device. The fourth factor is the tolerance to glitching. Glitching occurs when the content is not displayed at the rate it was authored causing the playback device to run out of data to display. In most cases any amount of start-up latency or glitching is intolerable, and it is therefore desirable to optimize the bit rate specified such that the start-up latency and the glitching are minimized or eliminated.\nWhatever bit rate is used, advertisements are generally not included into the media stream. Rather, the advertisements are inserted later, often targeted to particular people or regions. One common technique is to use multiple controls that compete for network bandwidth, which could delay the advertisement start-up and/or affect the quality of the advertisement rendered."}
{"text": "The present invention relates to protective cases for electronics. More specifically, the present invention relates to protective cases for electronic devices.\nElectronic devices such as cellular telephones, tablet computers, media players, global positioning system receivers, and personal data assistance devices are gaining widespread usage. These devices typically include a display, commonly referred to as a screen. In addition, these devices typically include a touchpad and one or more camera lenses. The screen, touchpad and camera lens are susceptible to various types of damage including scratches, nicks, impacts and chemical exposure. Though relatively fragile, these devices are increasing in costs. Accordingly, consumers desire to protect their products from accidental damage, and accordingly, there are now a wide variety of companies selling protective cases and films to help protect electronic devices.\nIn particular cellular telephones are increasing in popularity, including the need to protect cellular telephones. Modern cellular telephones typically incorporate a camera. These cameras have improved in quality to the point where many consumers have discarded or do not use a separate camera. Unfortunately, cellular telephone cameras still do not provide the picture quality of a traditional camera. This poor picture quality is primarily due to the comparatively poor lens quality within a cellular telephone's camera module.\nThus, there is a need to improve the picture quality of cellular telephone cameras. There is also a need to provide added protection to cellular telephone cameras."}
{"text": "The present invention relates to devices for operating an electrical switch in response to changes in sensed pressure in a cavity, commonly known as pressure switches, and particularly relates to pressure switches of the type having a pressure responsive member, such as a piston or flexible diaphragm defining a wall portion of the sensing cavity and movable in response to changes in pressure in the cavity. Pressure switches of the type employing a flexible diaphragm for sensing pressure are commonly employed in household washing machines for sensing pressure in a tube connected to the bottom of the washing machine receptacle or tub such that the sensed pressure is directly related to the height or depth of the water in the tub. If the washing machine tub is filled above a predetermined level, the hydrostatic pressure of the water fill is applied to the pressure switch diaphragm, and actuates an electrical switch for disabling the water fill valve or mechanism.\nIn the typical diaphragm-type pressure switch, the diaphragm is resiliently biased and preloaded by a spring; and, the preload is adjusted for calibration by an adjustment screw provided in a movable adjustment lever provided on the housing structure with the screw bearing against a movable spring reaction member such as a cap over one end of the preload spring. The movable adjustment lever is typically provided to enable user-adjustment of the pressure settings for switch actuation. In washing machine applications for the pressure switch, the user adjustment of the pressure settings corresponds to different desired levels of water in the machine wash tub.\nWhere an adjustment screw is used to alter the position of a reaction member attached to the preload spring, it is necessary to secure the adjusted position of the screw such as with an adhesive or mechanical crimping or staking. This construction has resulted in a costly and troublesome operation in mass production of pressure switches for household washing machines, and has also resulted in problems in field service where the means for securing the screw adjustment have failed, causing the pressure switch to change calibration. Thus, it has long been desired to provide an improved, more reliable and lower cost way of providing for calibration of a pressure switch."}
{"text": "Heretofore U.S. Pat. No. 4,446,158 has been obtained and discloses the use of a disposable individual coffee container and filter unit that includes a flat frame having an apertured center section with filter layers covering the top and bottom of the frame and used for retaining individual quantities of coffee or other material in compartments formed in the frame. This filter unit is positioned on the seat of a receptacle for such unit, which receptacle is adapted to be positioned on top of a cup or other container for receiving a quantity of liquid poured into the receptacle for flow through the filter unit for forming a desired beverage. In such a prior filter unit and apparatus for making a beverage, a by-pass opening can be provided therein whereby some of the liquid could flow through the apparatus without literally flowing through the beverage making material whereby the speed of making the product could be controlled depending upon the size of the by-pass slot. This partially controlled the strength of the beverage produced but a uniform or repeatable strength beverage would be produced in all instances by the apparatus. Inasmuch as it is well known that different people like different strengths of beverages, and particularly coffee, it would be very desirable if a type of coffee making apparatus could be provided wherein a varied strength product could be obtained. Possibly variable strength products could be obtained by the use of varied types of filter cartridges but this gets to be an expensive or cumbersome type of a manufacturing and marketing approach to the problem of making variable strength coffee and other beverages.\nIt is the general object of the present invention to provide a new and improved beverage making apparatus wherein a variable strength product can be obtained by use of the apparatus and by use of a standard type disposable, individual beverage making material container and filter unit.\nAnother object of the invention is to modify coffee making apparatus of a known type whereby a variable by-pass opening can be provided in the apparatus to aid in controlling the strength of the beverage produced and to aid in controlling the flow time of the beverage through the coffee making apparatus.\nYet another object of the invention is to provide a relatively inexpensive, uncomplicated modification in a beverage making apparatus for making small quantities of beverage such as coffee and wherein an adjustable by-pass opening is provided in the positioner ring or other part of the individual coffee container and filter units used.\nAnother object of the invention is to provide a movable member usually having a variable size flange or rib thereon for varying the flow through a by-pass opening in the beverage making apparatus, or to provide some type of apparatus with a plurality of different sizes of adjustable opening bypasses formed in the apparatus to enable a user to control the beverage making time and/or the strength of the beverage produced.\nStill another object of the invention is to provide an easily adjusted ring in a mounting plate used for positioning the individual disposable coffee carrier and filter units thereon and wherein a rotary adjustment device is provided on the plate in association with a by-pass slot or opening for controlling flow of liquid through the by-pass slot and which adjustment ring may include a variable height cut-off shoulder formed on the outer periphery of the ring for engaging the outlet of the by-pass slot and controlling liquid flow therethrough.\nThe foregoing and other objects and advantages of the invention will be made more apparent as the specification proceeds."}
{"text": "1. Field of the Invention\nThe present invention relates to a circuit structure for use in an electronic circuit such as a logical circuit operating at a high speed and a memory, and in particular to a circuit structure of a semiconductor integrated circuit.\n2. Related Background Art\nA semiconductor integrated circuit of the type described generally includes a logical circuit or a memory formed in a semiconductor chip. With a demand for digitalization in recent years, a request is increased for digital semiconductor integrated circuit operating at a high speed in a GHz band. The digital semiconductor integrated circuit generally includes a large number of transistors. In this event, each transistor is connected to a power supply for supplying electric energy and to a ground which discharges the supplied energy, thereby forming a transistor circuit.\nWhile such a transistor circuit operates as an on-off switch or a gate circuit, no signal is sent from the transistor even if the transistor is turned on unless the power source supplies the electric energy.\nConsidering the transistor gate circuit here, the gate circuit is instantly put into the on-state. If the on-resistance is low, such a phenomenon that the wiring is incapable of supplying electric charges occurs when the connected power supply attempts to supply electric energy (electric charges). It will be described later in detail why the wiring is incapable of supplying the electric energy. Likening the power/ground wiring to a transmission line and assuming its characteristic impedance is 50 Ω, if the on-resistance of the transistor gate circuit is lower than the characteristic impedance such as 15 Ω, the wiring is incapable of supplying the electric charges. Fortunately, the characteristic impedance of a signal line is 50 Ω or higher in many cases. Consequently, a problem of failure in supply can be avoided. It is necessary, however, to cope with an instant change for supplying charges to a transistor parasitic capacitor. At any rate, it is required to make a pipe for supplying charges thick, namely, to adopt a method and a structure for lowering the power/ground characteristic impedance. Therefore, the inventors of the present invention has proposed a method and a circuit structure for lowering the power/ground characteristic impedance in Japanese Unexamined Patent Publication (JP-A) No. 2000-174505.\nFurther, if a switching operation of the transistor circuit becomes faster, an inductance against a sharp increase of current can not be ignored with the above-mentioned problem. Even if the power/ground characteristic impedance is lower than that of the signal line, it cannot catch up with rapid opening of a transistor gate circuit and a parasitic capacitance charge due to the parasitic inductance of the power/ground wiring. Therefore, a method and a structure for lowering the parasitic inductance are necessary.\nMoreover, signal energy (charge quantity) determining a reception is insufficient until charging all the capacitors in the receiving transistor gate is completed with the current controlled by on-resistance. This causes an operation delay of the receiving transistor itself. As a result, current of the power supply continues to flow during the time.\nIn this manner, a state of disabled smooth control of switching in the transistor gate circuit becomes prominent in a digital circuit of the GHz band due to rate controlling caused by a supply capacity of the power/ground."}
{"text": "This invention relates to an anti-fouling overcoating composition for watercraft and underwater portions of marine structures precoated with a hard-surface leaching-type anti-fouling paint and to a method for reducing fouling and, secondarily, reducing or minimizing drag created in fluid passing along the surfaces of underwater structures, whether movable or stationary.\nMarine fouling organisms, such as crustaceans (barnacles, Conchodermae), mollusks (boring worms and boring clams), annelids (tubeworms, Serpulidae), tunicates, algae, coelenterates, acorn shells (Balonidae), goose mussels (Lepodoids), and moss (Hydroidae) grow and accumulate on surfaces in an underwater environment. Wooden structures are particularly susceptible to attack by pileworms or shipworms (Teredinidae), limnoria, martesia, sphaeroma and the like. Metal water intake pipes for cooling towers in thermoelectric power stations, petrochemical factories of the like are commonly infested with and fouled by adherence of Blue mussel, sessile acorn barnacles (Balanus), oysters, moss (Bryozoa), Hydrozoa and the like.\nApproaches to prevention of a marine fouling have included the use of silicone resins, as disclosed by Mueller et al (U.S. Pat. No. 3,702,778); the use of asphaltic bitumen compositions as recited by Rudd (U.S. Pat. No. 2,489,228) and compositions containing antifoulant in a hydrophilic polymer derived from a hydroxyalky acrylic polymer (Shepherd et al U.S. Pat. Nos. 3,990,381; 3,896,753; and 3,575,123).\nVassileff (U.S. Pat. No. 3,786,113) indicates that polymer compositions having selective permeability for gas and water vapors are obtained by mixing a polyacrylic resin and polyethylenimine, in 5:1 to 2:1 weight ratio, and cross-linking the resulting mixture.\nMerlino et al (U.S. Pat. No. 3,702,349) employ the reaction product of a polyaziridinyl adduct and a coreactant containing carboxy or other functional groups in an organic solvent-based coating material.\nWalus (U.S. Pat. No. 3,821,145) describes an aqueous coating composition of an acrylic polymer containing hydroxyacrylate or acrylic acid units and a water-dispersible methylolated melamine resin.\nHurwitz (U.S. Pat. No. 2,954,358) discloses aqueous compositions comprising a water-soluble or readily water-dispersible polyepoxide and a water-insoluble polymer containing 1-10 mole percent of acrylic or methacrylic acid monomer units. The use of amine, including polyalkylenepolyamines, in small amounts as a catalyst for the cross-linking reaction between the epoxide function and the acid function is set forth.\nBlank (U.S. Pat. No. 3,959,237) discloses the use of a film-forming copolymer containing 20-60 parts of acrylic or methacrylic acid, at least partially neutralized by multivalent cations, as a carrier for agricultural fine chemicals. Ethylenimine is disclosed as neutralizing agent.\n\"Fouling\" can be defined as assemblages of marine animals and plants which grow on watercraft and underwater marine structures, as well as on rocks, stones and other natural objects. Marine fouling is a result of growth of marine animals and plants on the exposed surfaces of man-made marine structures, including bottoms and hulls of ships, pier supports, buoys, water-intake pipes, fishing nets, rowboats, water skis, ocean liners, tankers and other cargo ships, submarines, pilings, bridge substructures and the like. Marine fouling is an omnipresent problem, occurring whether the underwater surface is made from wood, metal, plastic, fiberglass, concrete or other structural material.\nFouling is economically harmful to the shipping industry because of the accompanying decrease in the speed of watercraft, resulting in higher power consumption and operating costs to maintain schedules, and because of the need to bring ships into drydock to remove the fouling and apply fresh paint. Furthermore, fouling can produce physical damage to watercraft or stationary underwater structures, resulting in shortened useful life and the need for frequent replacement. In the case of fishing nets, fouling causes an increase in the fluid resistance to water flowing therethrough, so that fish growing in a limited area, as is typical of areas of cultivation, are adversely affected by a poor supply of oxygen in the water.\nAttempts to control fouling by marine organisms generally have been based on the use of chemicals which are toxic to a specific organism or group of organisms. Thus, creosote has been used for the impregnation of wooden surfaces. However, creosote is unsatisfactory for painted surfaces, for the reason that creosote bleeds through the paint.\nA variety of anti-fouling agents have been developed which are toxic to deleterious marine organisms. Typical of these agents are cuprous oxide, mercury oxide or the like; organocopper salts, e.g., copper naphthenate, copper oleate, etc.; organotin compounds, e.g., bis-tributyltin oxide, triphenyltin bromide, dibutylethyltin bromide; 1,2,3-trichloro-4,6-dinitrobenzene; dichlorodiphenyldichloroethane; nitrodiphenyl ether derivatives; organolead compounds, e.g., triphenyllead stearate, triphenyl lead chloride; 10,10'-oxybisphenoxazine (SA-6 546); hexachlorophene and tetrachloroisophthalonitrile, alone or in combination. However, most of these materials are highly toxic to other organisms and tend to leach out of the coating in which they are applied so that the coating composition becomes ineffective in preventing fouling after a period of time. Often, compositions containing these materials are difficult to apply to marine structures for lack of adequate adhesion or tend to peel from the structure after an unsatisfactorily short interval.\nIt is therefore apparent that there is a continuing need for anti-fouling systems which adhere tenaciously to the surface treated therewith, operate efficiently as anti-foulants over a prolonged period of time and are not highly toxic to desirable forms of marine life."}
{"text": "The present disclosure relates generally to printheads.\nInkjet printing creates images by propelling ink droplets onto a medium. An inkjet print head includes an array or a matrix of ink nozzles, with each nozzle selectively ejecting ink droplets. The number of operating nozzles and the drop volume establish the ink flow from an ink reservoir or supply, which may be an intermediary ink tank placed in close proximity to the print head or a remote ink tank. When printing average density images, the print head tends to consume steady amounts of ink. However, sudden changes in ink consumption often occur at the beginning and the end of the printing process. The energy used during an ink firing event may create motion of the ink in the firing chamber and ink delivery system, which may cause fluidic interaction between neighboring ink channels. When the interactions are large enough, crosstalk may occur, where the firing event of one channel may cause a disturbance in neighboring channel firing events."}
{"text": "The invention relates to computer data storage devices. More specifically, the invention relates to a tape cartridge for a computer tape drive.\nMajor causes of data loss in computers include accidental file deletion, hard-drive malfunctions, viruses and system failures, power surges, theft of computers and natural disasters such as fires. Installing and uninstalling software can also cause data loss.\nLoss of data stored in a computer can result in a significant loss of time and money. For a business, any of these causes might result in the loss of client lists, product information, payroll, personnel files, tax records and other vital information. The consequences can be disastrous.\nTape technology offers efficient and cost-effective protection against data loss. Tape drives offer a range of media that allow data on hard drives and servers to be backed up. A snapshot of a full system, including programs and data files, can be captured on a single tape cartridge. Backing up a computer on a regular basis will ensure that vital data is not lost.\nA typical tape cartridge includes magnetic tape and one or two reels. A Travan cartridge is a well known type of two-reel cartridge. A \"3480/3490\" cartridge is a well known type of single reel cartridge.\nThe tape must not become damaged or degraded. If the tape becomes degraded or damaged, valuable data might be lost forever.\nThe tape can become damaged or degraded by debris inside the tape cartridge. During storage or use of the tape cartridge, debris such as particles of dust and dirt can enter inside the tape cartridge. Such debris can damage the tape and interfere with the read/write process. During use of the tape cartridge, oxide particles can fall off the tape. These wear particles can also damage the tape and interfere with the read/write process.\nProblems with debris will become even greater as storage capacity of the tape is further increased to accommodate the ever-increasing storage capacity of hard drives. The storage capacity of the tape may be increased by increasing bit density. As the bit density is increased, however, the tape becomes more sensitive to wear particles and the other debris inside the cartridge.\nThere is presently a need to protect magnetic tape against debris. This need will become greater as tape bit density is further increased."}
{"text": "As a result of increased environmental awareness, water conservation has become an important issue especially in the Western part of United States, densely populated areas elsewhere within the United States and in other countries or regions that have arid climates.\nWith regard to domestic water consumption in this country, the ultimate use is generally for sanitary, culinary, drinking, washing or bathing purposes. It has been found that careless or willful waste occurs during the performance of certain of these tasks such as washing dishes, preparing food, brushing teeth, etc. wherein the user permits the water to flow continuously from the faucet rather than to temporarily shut-off the flow, when water is not needed, during these procedures. The primary reason for not interrupting the water flow is generally because the user's hands are occupied or because of the effort required to again achieve the same water temperature and/or pressure.\nBy way of example, water consumption during the process of washing and rinsing a typical load of dishes for four place settings, consumes 7.5 gallons of water or approximately one cubic foot of water. During the wash/rinse period, the water is idling for about fifty (50%) percent of the time, thus wasting 0.5 cubic feet.\nAside from the environmental issues, another motivation for reducing consumption is economic in nature in that there is a tendency for the rates charged for metered water to generally increase rather than decrease. Furthermore, the conservation of hot water minimizes the fuel needed to heat the water and thus achieves additional cost savings.\nAn apparatus to control a water faucet valve without requiring hands-on operation is shown in U.S. Pat. No. 2,270,239. A shortcoming of that device is that the moveable components of the mechanical linkage are subjected to friction and wear resulting in slack and inefficiency. A further problem with that device is that the installation requires a plumber or other skilled workman and it is not particularly adapted as a retro-fit assembly.\nAnother remote control faucet valve device, as shown in U.S. Pat. No. 4,052,035, utilizes a flexible conduit for communicating between a foot control member and a valve member. A disadvantage of that arrangement is that the conduit is draped over the sink and the kitchen cabinet when in use and thus presents a physical impediment or hinderance which interferes with access to the sink and underlying cabinet. The conduit can also further become entangled with kitchen appliances. Additionally, the hydraulic system has a restricted range of control functions.\nWith regard to wireless control systems utilizing a radio link, these devices have been applied to the operation of motorized toy vehicles and model airplanes. They have also been applied to remote keyless entry systems for locking or unlocking the doors and trunk lid of automobiles. The utilization of radio controls for home automation has been rather limited with the most common usage being for electronic garage door openers.\nFluid handling apparatus employing radio signal control has been applied to irrigation systems as discussed in U.S. Pat. Nos. 3,726,477 and 4,838,310. The systems of the aforementioned patents, however, are not adapted for domestic use and do not include an electronic valve structure or a tactile foot-operated sending unit as in this invention."}
{"text": "The invention generally relates to automotive safety restraint systems and more particularly to a pretensioning system capable of simultaneously removing the slack in a plurality of seat belts, secured about respective occupants, utilizing a single pretensioning or belt-tightening device.\nPretensioners or belt tighteners, as they are also known in the art, are used to tighten a portion of a seat belt system, thereby removing slack about the occupant, during an accident. This action is done prior to the time the occupant begins to move forward as a result of the deceleration of the vehicle.\nAs a typical accident progresses, the occupant will remain fixed to the seat and begin to move forward at about 20 milliseconds after the initiation of the crash. Belt tighteners must react and eliminate the slack in the seat belt system, about the occupant, prior to the occupant's moving forward. Typically the pretensioning system is designed to react quickly and eliminate the slack in about 10 milliseconds upon receipt of a control signal, which is generated by one or more crash sensors and associated electronic control systems. There exists in the art two basic types of seat belt pretensioning systems. The first is associated with the spool of a retractor and the shoulder belt of a three-point seat belt system. During the accident the spool is caused, by various means, to reverse-wind or retract thereby removing the slack, typically in the shoulder belt portion of the seat belt system, to thereby hold the occupant in place. The other type of pretensioning system is called a buckle pretensioner. The buckle pretensioner is typically connected by a length of seat belt webbing or by a cable to the seat belt buckle. These buckle pretensioning systems (as well as retractor systems) include a pyrotechnic element which is fired, causing the buckle to be pulled downwardly, thereby reducing slack about the occupant. The prior art generally shows one pyrotechnic element associated with each shoulder belt (for the retractor pretensioner) or with each buckle (for a buckle pretensioner).\nIt is an object of the present invention to provide a single pretensioning device that is capable of eliminating the slack in a plurality of seat belt systems.\nAccordingly the invention comprises a system for eliminating residual seat belt slack about a plurality of occupants of a vehicle, the system comprising: a plurality of seat belt assemblies securable about each respective occupant, each seat belt assembly including a lap belt, a tongue, and a buckle. One end of the lap belt is secured to a support member and an opposite end of the lap belt is operatively joined to the tongue. The buckle matingly locks with the tongue and the buckle is secured to one end of a connecting member, the connecting member having an opposite second end. A belt tightening means, secured to the second ends of each connecting member, generally simultaneously tightens each seat belt assembly about a respective occupant. Each belt assembly may also include a shoulder belt.\nMany other objects and purposes of the invention will be clear from the following detailed description of the drawings."}
{"text": "1. Field of the Invention\nThe present invention relates to removal of deposits from the interior of a tobacco drying apparatus. More particularly, the present invention relates to the removal of tobacco deposits in a high humidity drying apparatus with the use of recirculating heated air.\n2. Discussion of the Prior Art\nIn the manufacture of cigarettes and like articles, it is the usual practice to reduce tobacco, the term being used herein to include both lamina and stems, to a particle size appropriate for manufacturing cigarettes. The moisture content of the tobacco is generally increased prior to this size reduction processing in order to minimize tobacco breakup and provide a material of uniform particle size. Furthermore, in order to process the treated tobacco in the manufacturing of cigarette rods, it is necessary to reduce the moisture content of the tobacco to a level below that which the tobacco is at after treatment by casings, flavorings and other high moisture containing additives. The actual drying process has a direct impact upon the quality of tobacco utilized during cigarette manufacturing because of the effect the drying process has upon the tobacco material itself.\nAdditionally, drying of tobacco after the addition of flavorings and casings has a direct impact upon the quality of the tobacco. If, during drying, the tobacco is subjected to rigorous agitation or contact with stationary surfaces, the tobacco material can be damaged by breakup thus decreasing the filling capacity of the tobacco. This unwanted result is also achieved when drying under low humidity conditions. It is therefore desirable to dry the moist tobacco under high humidity drying conditions.\nThe high humidity pneumatically conveying drying apparatus used to dry cut tobacco typically requires cleaning due to the interior buildup of deposits of flavorings and other materials added to the tobacco as well as portions of the tobacco material itself. The buildup of material in the dryer is most prevalent in areas where the direction of the pneumatically conveyed tobacco is changed. These `corners` or `elbows` usually have the highest incidence of material buildup on their interior surfaces and this buildup can cause hygiene problems as well as reduce the smooth flow of material coming into contact with the built up residue.\nPreviously, to remove the buildup of material in tobacco dryers, whether they be high humidity dryers or typical in-line dryers, the dryer had to be shut down and cleaned out by flushing with a high volume of hot water or other solution at relatively high pressure. Usually, this cleaning process was performed every 15 to 17 drying operations. Thus, this cleaning had the undesirable effect of lengthy down time for the drying apparatus while the dryer was cleaned. Additionally, if a cleaning solution was used, water had to be flushed through the dryer to remove all traces of the cleaning solution so that subsequently dried tobacco would not be contaminated with the cleaning solution. The excessive cleaning time severely impacts the manufacturing of tobacco products. Moreover, liquid cleaning does not contact all interior drier surfaces.\nPrior art attempts to clean interior surfaces and passageways of equipment is shown in U.S. Pat. No. 3,888,302 wherein a heat exchanger for coke ovens and interior passageways therefor require general cleaning on a periodic basis. A reactive material, such as anhydrous hydrochloric acid, is introduced into the passageways to produce reaction products which are volatile when removed in a gas stream and thereby discharged. However, the ability to introduce such reactive material into the present system is extremely restricted due to the necessary requirement of completely removing the reactive material prior to further tobacco processing."}
{"text": "1. Field of the Invention\nThis invention relates to reduced dimension MOS field effect transistors and to the formation of field effect transistors having narrow gate electrodes and reduced dimension source/drain structures.\n2. Description of the Related Art\nField effect transistors, generally referred to as FETs or MOSFETs, are the most common devices in modern integrated circuit devices. One conventional configuration of a FET is illustrated in FIG. 1. Field isolation regions 12 are formed by the local oxidation of silicon (LOCOS) method at the surface of a substrate 10, defining the active device regions and providing lateral isolation between adjacent devices formed in and on the surface of the substrate 10. A gate oxide layer for the FET covers the active device regions of the substrate 10 and a gate electrode 16 of doped polysilicon is formed on the gate oxide layer 14. Oxide spacer structures 18 are generally provided on either side of the gate electrode 16. The inner edges of source/drain regions 20 define a channel region at the surface of the substrate, with a source/drain region extending from either side of the gate electrode 16 to the field isolation regions 12. Frequently, the source/drain regions 20 have a lightly doped drain (LDD) structure in which an inner, more lightly doped portion of the source/drain region is aligned with the edge of the gate electrode 16, and a more heavily doped portion of the source/drain region 20 is aligned with the oxide spacer structure 18. A P-pocket or halo implant 22 extends slightly below and toward the FET channel from each of the LDD source/drain regions. The P-pocket or halo implant reduces the short channel effects in the FIG. 1 FET.\nGenerally, the FET structure shown in FIG. 1 is prepared by first forming a field isolation mask on the surface of the substrate 10, with openings in the mask exposing the substrate over regions where the field isolation structures will be formed. The field isolation structures are then formed either using a local oxidation of silicon (LOCOS) process, as illustrated, or a shallow trench isolation method. The field isolation mask is then stripped and various implantations may be made into the active regions of substrate 10 to adjust the doping profile of the substrate within the active regions. A gate oxide layer 14 is then grown on the cleaned surfaces of the active regions of the substrate 10. Polysilicon is blanket deposited by a low pressure chemical vapor deposition (LPCVD) technique over the gate oxide layer and the field isolation regions. The polysilicon layer is doped, typically by ion implantation, and photolithography is used to define gate electrodes 16 over the active regions. The source/drain regions 20 are formed in a two-stage implantation process. A first ion implantation is made using the gate electrode and the field isolation regions to mask the substrate, forming the more lightly doped portions of the LDD source/drain regions 20. A layer of CVD oxide is then deposited over the gate electrodes and over the surface of the device. An etch back process forms spacer structures 18 from the CVD oxide layer on either side of the gate electrode 16. A second ion implantation is performed to form a more heavily doped region aligned to the oxide spacer structures 18 and completing the source/drain regions 20. In the illustrated NMOS device, the source/drain regions 20 may be doped with any N-type dopant or combinations of different N-type dopants might be used to achieve different diffusion profiles. Finally, the P-pocket region 22 is formed by angular ion implantation of boron ions with the implantation made so that the ion flux is at an angle of approximately 30.degree. from normal (perpendicular) to the substrate.\nImprovements in device density and reductions in the cost of manufacturing integrated circuits are closely linked to reducing the size of devices within the integrated circuits. The width of gate electrode 16, as well as the sizes of other device structures, are determined by conventional lithography processes. Shrinking the size of the FET shown in FIG. 1 generally cannot proceed beyond the resolution and alignment limitations of the particular process technology used in forming the FIG. 1 device. Thus, the width of the gate electrode 16 is typically designed to be a width d equal to the design rule for the particular process used in making the gate electrode. Further reductions in the size of the gate electrode are desirable to decrease device size and to improve the density of the integrated circuit. Adoption of higher resolution lithography techniques, which could facilitate forming smaller gate electrodes, is very expensive, and may only be economically justified for very high volume manufacturing. Smaller volume manufacturing operations and specialty or low-profit margin circuits might not cost effectively implement such high cost processes. Thus, even with the introduction of higher resolution lithography techniques, it may be difficult to further reduce the size of the FIG. 1 FET.\nAnother disadvantage of the device illustrated in FIG. 1 and of the method of making the device is that the source/drain electrodes require a significant implantation to ensure that their resistance is sufficiently low to provide good device performance. The required high level of ion implantation causes a variety of problems. For example, the heavy ion implantation dosage renders the substrate amorphous where the source/drain regions are to be formed. Recrystallization of the substrate in the source/drain regions is then performed in an annealing process which can produce defects in the recrystallized material, or which can lead to excessive levels of diffusion from the source/drain regions. Excessive diffusion from the source/drain regions can make the channel region beneath the gate electrode 16 narrower than is desired, compromising device performance. A further disadvantage with the FIG. 1 structure is that the P-type halo implant 22 overlaps with a considerable portion of the source/drain region 20. This extended overlap is undesirable for a variety of reasons, including the size of the P/N junction formed between the halo implant 22 and the source/drain region 20, which impairs the FET's performance by reducing switching speeds."}
{"text": "Freeze-drying processes are frequently used to preserve microorganisms. Indeed, the low operating temperatures of freeze-drying processes lead to minimal damage in microorganisms. Moreover, it is known in the art that low storage temperatures can be an essential element in maintaining the viability and stability of such freeze-dried cells. It is also known in the art, that cells which are freeze-dried in the presence of protective agents are better suited to maintain their viability and stability than cells which are freeze-dried without the addition of said protective agents. For example, U.S. Pat. No. 6,919,172 (Wisconsin Alumni Research Foundation) relates to the use of a protective mixture (i.e. sucrose and potassium phosphates) mixed with biological material in order to form a preservation medium, prior to the biological medium undergoing at least one preservation process (eg. freezing, freeze-drying, vacuum-drying or spray-drying). After being subjected to a preservation process, the biological material is capable of remaining stable when stored in dry conditions at room temperature. However, the probiotic microorganisms of U.S. Pat. No. 6,919,172 which have been preserved to provide increased stability during shipping and storage in sealed containers cannot usually withstand the simultaneous high temperature and/or high humidity conditions which are typically encountered in long-term storage and secondary processes used in feed products, food products (e.g. cereals, nutritional bars, recombined cheese, infant formulas), consumer healthcare products (including dietary supplements, nutraceuticals, over-the-counter products) as well as agri-products. Indeed, the viable microoganisms are often inactivated by heat and/or moisture and/or other processing conditions or storage conditions, thereby decreasing their viability in the final product over time, during storage and/or during consumption, in such a way that the desired effect of the final product containing the microorganisms is eventually lost.\nIn order to overcome this problem, numerous coating processes for microorganisms have been developed to protect them against undesirable elements such as high temperature and high humidity levels. These coating processes are often based on coating materials such as fats or low permeability polymers such as shellac and hydroxypropylmethyl cellulose (HPMC). For example, U.S. Pat. No. 7,157,258 (Lallemand) deals with coating of dehydrated microorganisms by using hydrophobic barriers such as fat and waxes with melting points greater than 60° C., while EP 1753440 B1 (Merck), U.S. Pat. No. 6,365,148 (Il Yang Pharm Co. Ltd.) and EP 0778778 B1 (University New South Wales) use polymeric barriers such as HPMC or starch. These moisture barriers are intended to provide protection of the dehydrated microorganisms by preventing or delaying penetration of water into the coating, thereby avoiding diffusion of moisture to the underlaying dehydrated microorganism. If the coating allows moisture to pass, then the dehydrated microorganism may absorb moisture, thereby altering its viability.\nHowever, the Applicant has prepared dehydrated microorganisms coated with commonly used compounds such as fats, shellac or HPMC and noticed that the viability of such microorganisms containing products was not sufficiently maintained. Indeed, as can be seen in the results section of the present patent application (especially Example 4, Table 5 and Example 7, Table 8), the materials used in the prior art to coat the microorganisms do not enhance the thermostability of the microorganisms nor provide enhanced viability over a longer period of time in medium or high humidity environments when compared to non-coated dehydrated microorganisms.\nConsequently, despite the teachings above, there still exists a need in the art for coated dehydrated microorganisms, whereby the microorganisms are able to retain most of their activity during and after manufacturing and are able to resist the damaging conditions encountered for example in secondary processes used for the preparation of feed/food products, consumer healthcare products or agri-products or during long term storage of these products, wherein said processes/long term storage often involve a combination of high temperature and/or high moisture, thereby responding to a widely sought-after and currently unmet demand for dehydrated microorganisms with improved stability."}
{"text": "Fluorescent reagents enable life science research in many fields, including biological, biomedical, genetic, fermentation, aquaculture, agriculture, forensic and environmental applications. Fluorescent probes and stains identify biopolymers and detect particular biological components within and outside cells. A common example is the use of fluorescent-labelled antibodies to detect cell-surface receptors. Another example is the widespread use of gel electrophoresis for characterizing nucleic acids, one limitation of which is the sensitivity of the methods used to detect the nucleic acid bands.\nDetection of biological analytes utilizing fluorescent labels eliminates the need for radioactive labels, thereby enhancing safety and diminishing the adverse environmental impact and costs associated with radioactive waste disposal. Examples of methods utilizing fluorescent detection methods include automated DNA sequencing, oligonucleotide probe methods, detection of polymerase-chain-reaction products, immunoassays, and the like. In the life and medical sciences, researchers and technicians often need to detect proteins, antigens, and other ligands on the surface of cells. Receptor based assays utilize labelled molecules, e.g. fluorescent labelled peptides, proteins, and antibodies to detect expressed proteins and other ligands.\nDyes that are generally applicable for staining or labelling biopolymers across a broad range of applications preferably have the following properties: (i) the dye-biopolymer conjugate or complex should produce a very high signal with low background so that small quantities of biopolymers can be sensitively detected in both cell-free and cell-based assays; and (ii) the conjugate or complex should be photostable so that the fluorescent signal may be observed, monitored and recorded without significant photo bleaching. For applications involving cell surface binding of dye-peptide or dye-antibody conjugates to membranes or cell surfaces, especially live cells, the dyes preferably (iii) have good water-solubility to achieve effective conjugate concentration and detection sensitivity and (iv) are non-toxic to living cells so as not to disrupt the normal metabolic processes of the cells or cause premature cell death.\nWhile many dye compounds have found use as nucleic acid stains, most available dyes fluoresce in the green region of the visible spectrum. Green lasers are more expensive than red lasers and give higher background signals in live cell assays due to autofluorescence of cellular components and assay equipment. These higher background signals decrease the sensitivity of the assay. Moreover, many cellular components absorb green light, further reducing the sensitivity of the assay. Thus, sensitive dyes that are photostable, have excitation and emission maxima in the red region of the visible spectrum and that are water-soluble are highly desirable."}
{"text": "Mobile phones now have functionality and applications that provide a wide variety of communication modes. For example, a single device can now support conventional phone calls, Voice-over-Internet-Protocol (VoIP) calls, video calls, and the like. However, such functionality has not been particularly well integrated, and new users may not even be aware that such functionality is available to them.\nBecause users can face hurdles when attempting to take advantage of different communication modes, there remains room for improvement."}
{"text": "The discussion below is merely provided for general background information and is not intended to be used as an aid in determining the scope of the claimed subject matter.\nAn aspect of the present invention relates to a method for the preparation of a highly concentrated liquid lactic acid composition and a composition obtainable by said method.\nLactic acid as well as its salts and esters have long been used as food additive and in various chemical and pharmaceutical applications. More recently, lactic acid has been used in the making of biodegradable polymers both as a replacement for present plastic materials as well as various new uses where biodegradability is needed or desired such as for medical implants and slow release drugs. Accordingly there is an ever-increasing demand for improved and economically viable lactic acid production processes.\nLactic acid is commonly marketed as a dilute or concentrated aqueous solution, whereby concentrated lactic acid solutions generally have a lactic acid concentration of about 90% (w/w). The marketing of lactic acid compositions having even higher concentrations is still desirable, e.g. to decrease packaging and transport costs and to improve ease of use in all kinds of technical and food applications. However, higher concentrations are not feasible because crystallisation of lactic acid in such highly concentrated compositions easily occurs at an ambient temperature, which provides problems with handling the lactic acid composition.\nIt is desirable to provide a highly concentrated liquid lactic acid composition that does not crystallise at an ambient temperature."}
{"text": "Intra-aortic and intra-ventricular cardiac assist devices are well known in the art. These devices are generally used to reduce the heart's work load after insult or surgery. They may also be used to increase blood flow from the left ventricle of the heart into the aorta in cases of insufficient cardiac output due, for example, to acute or chronic heart ailments or to interference with normal cardiac function during surgery.\nOne of the best-known and most widely-used intra-aortic pump systems is the Intra-Aortic Balloon Pump (IABP), comprising a catheter, having an inflatable balloon at its distal end, which is inserted through an artery into the aorta. The balloon is then alternately inflated and deflated by an external pump drive, so as to alternately block and unblock blood flow through the aorta, in synchrony with the beating of the heart, in order to assist the left ventricle in propelling blood into the arterial system. The IABP, however, provides only limited augmentation of the heart's natural, unassisted output, and is not adequate for overcoming heart failure.\nU.S. Pat. No. 4,014,317, which is incorporated herein by reference, describes a cardiocirculatory assist cannula with a balloon pump and cardiac pacing electrode. The cannula is inserted percutaneously through the aorta so that its distal end is inside the left ventricle of the heart. During systole, inlet valves on the cannula inside the left ventricle open, and the contraction of the ventricle forces blood to flow into the cannula. Then, during diastole, the blood flows out, into the aorta, through one or more outlet valves along the cannula downstream from the inlet valve. A gas-filled balloon, similar in function to the IABP described above, is connected to the cannula downstream of the outlet valves. The balloon is typically inflated during diastole and deflated during systole, to assist in perfusion of the coronary arteries. The cannula has a small stroke volume, however, and relies on the contractile force of the heart to pump the blood. It is therefore of limited usefulness in augmenting the blood output of a weakened or failing heart.\nU.S. Pat. No. 4,906,229, which is also incorporated herein by reference, describes a high-frequency transvalvular axisymmetric blood pump. The pump includes a small internal volume, which may be alternately expanded and reduced by pneumatic or hydraulic pressure which is exerted via a flexible membrane radially surrounding the volume. The volume has intake and outlet ends, with one-way axial valves at both of the ends, so that blood can flow only from the heart into the aorta. The pump is connected via the one-way intake valve to a cannula, which is inserted into the left ventricle of the heart through the aortic valve. When the internal volume is expanded, blood flows into the pump from the ventricle. The volume is then reduced, and the blood is ejected into the aorta through the outlet end. This pump is designed to operate at a frequency of 600 to 1,000 cycles per minute. Since the stroke volume of the pump is typically only about 3-5 cc, these high cycle rates are needed in order to provide adequate perfusion.\nIn the Hemopump Cardiac Assist System, distributed by Johnson & Johnson Interventional Systems, a cannula containing a special, miniature rotor pump mechanism is inserted into the aorta. The pump is driven by a drive unit outside the body, to pump blood continuously from the aorta into the rest of the arterial system, thereby supplementing the heart's natural output rate. A system of this type is similarly described in U.S. Pat. No. 5,092,844, which is incorporated herein by reference. While continuous-flow devices are useful for short-term augmentation of cardiac output, it is believed that pulsatile pumps provide more effective long-term support, since they approximate more closely the natural pump action of the heart."}
{"text": "1. Field of the Invention\nThe present invention relates to a paper feeding apparatus for an image forming machine such as a copying machine or a laser printer.\n2. Description of the Prior Art\nIn a conventional copying machine, in a place where the lateral width of the copying machine body is restricted for installation, in order to insert or to take out a paper feeding hopper 3 such as a paper feeding cassette from the front side (right side in FIG. 1), a so-called front loading cassette type has been employed, wherein, as shown in FIG. 1, the paper feeding hopper 3 such as the paper feeding cassette is inserted and taken out freely from the axial direction 61 of a roller shaft 2 of a paper feed roller 2 relative to the paper feed roller 1 supported rotatably by the image forming machine body. Papers 4 are biased upward by means of springs 22.\nIn the prior art aforementioned, when the paper feeding hopper 3 placed on a drawer is used as it is, as shown in FIG. 1, the papers 4 are moved in parallel to the axial direction 61 of the paper feed roller 1 and caught by the paper feed roller 1 as shown by reference numeral 4a, thus causing a slack in the paper feeding hopper 3 and jamming during feeding of the paper.\nAlso, such a problem is encountered that a paper feeding hopper for a drawing mechanism has to be newly designed, resulting in high cost."}
{"text": "The present invention relates to the field of electrical connectors, especially for telephone communication equipment, and more particularly to environmentally protected electrical connections. Most particularly, in one embodiment the present invention provides a method and apparatus for protecting telephone jacks from moisture and corrosion, such as often found in coastal regions, islands, and the like.\nTelephone line connections at subscriber locations are commonly made with the RJ11 type of plug and socket connector. These connectors are exemplary of electrical connections susceptible to failure from oxidation, corrosion, humidity, salt, and the like, especially in the presence of a live voltage on the conductors within the connector.\nFor example, it is sometimes difficult to establish and maintain an adequate environmental seal in a removable male RJ11 plug, particularly when wires lead from the male RJ11 plug. Accordingly, moisture and other environmental contaminants are allowed to enter such plugs, sometimes resulting in corrosion and/or failure of the connection of the tip and ring connections in the socket/plug combination. RJ11 sockets are likewise subject to moisture contamination and corrosion, as well as being subject to dust buildup. In hot, humid environments such as in Florida and along the Gulf Coast of Texas, failure can occur within several months of installation. Servicing these failures is costly for the consumer or the telephone company.\nSometimes problems have also arisen in connection with test ports for customer telecommunications equipment such as remote terminals at customer facilities, described in the parent application, and the like. It is often desirable to provide an RJ11 connector of the type well known to those of skill in the art, or other such connector, at an external location at subscriber facilities such as a junction box leading to a house or a remote terminal of the type described above. Previously, such access is provided by installing a female RJ11 socket at such locations which is normally connected to a male RJ11 plug. The tip and ring wires (among other wires in some cases) lead from the female RJ11 socket, and connect to tip and ring connections in the male RJ11 plug, thereafter leading into the subscriber facility. When it is desired to connect test equipment to the RJ11 female socket, the plug is removed, and another male RJ11 is inserted into the female socket, thereby providing tip and ring connections for the test equipment. Even though the equipment may be contained in a protective housing, such arrangements are sometimes subject to much of the same moisture/corrosion degradation.\nIt would therefore be desirable to provide an improved method and associated apparatus for protecting plug and socket electrical connectors from the environment. In particular, an environmentally resistant RJ11 plug and socket apparatus would be especially desired."}
{"text": "The generation of immunity to infectious organisms is a powerful tool in disease control. Those antigens that induce immunity to infection are known as immunogens. The protective antibody they induce may collaborate with other natural defenses to inhibit the infective process, or they may neutralize harmful products of the infective organism such as toxins.\nAn effective means of enhancing the antibody response is the use of an adjuvant. Thus, an adjuvant is included in a vaccine as an additive or vehicle to enhance the response to the antigen. An adjuvant may function by different mechanisms, including (1) trapping the antigen in the body to cause a slow release, (2) attracting cells of the immune system to the injection site, (3) stimulating cells of the immune system to proliferate and to become activated, and (4) improving antigen dispersion in the recipient's body.\nA number of agents with diverse chemical properties have been used as adjuvants, including water-in-oil and oil-in-water emulsions, mineral salts, polynucleotides and natural substances. One adjuvant, known under the trademark AMPHIGEN™, is described in U.S. Pat. No. 5,084,269. AMPHIGEN™ adjuvant consists of de-oiled lecithin dissolved in an oil, usually light liquid paraffin. In vaccine preparations AMPHIGEN™ is dispersed in an aqueous solution or suspension of the immunizing antigen as an oil-in-water emulsion.\nProblems were observed when using an AMPHIGEN™ adjuvant according to U.S. Pat. No. 5,084,269, above. For example, the lecithin in the AMPHIGEN™ does not suffice to produce a stable emulsion of the oil, thus leading to a pool or depot of oil in the injected tissues. Mineral oil can not be metabolized or removed by the animal. As a result, the oil becomes a source of severe chronic inflammation and scarring. Emulsifying the AMPHIGEN™ directly in the antigenic preparation carries the risk of damaging the antigen. Also, if the desired emulsion fails to form, the valuable antigen must be discarded.\nAn adjuvant useful in vaccines for animals, including humans, that is effective and solves the above problems would therefore be highly desirable."}
{"text": "Various types of automatic beverage dispensers are known. Generally described, these dispensers hold a measure of dispersible materials in a container of some sort. Hot water typically is added to the materials so as to mix the beverage. One drawback with these known beverage dispensers is that the elements of the dispenser that come in contact with the dispersible materials must be cleaned and/or sanitized on a periodic basis. Further, dispersible materials generally require a significant amount of work to properly mix the beverage. As a result, the beverage dispenser as a whole may be somewhat slow between beverage cycles.\nThere is a desire, therefore, for a beverage dispenser and associated components that mixes a beverage with a relatively quick cycle time. The beverage dispenser preferably should be relatively inexpensive and easy to use while consistently producing a high quality beverage. Likewise, the beverage dispenser preferably should be easily adaptable for different types and amounts of dispersible materials and other ingredients."}
{"text": "The present invention relates to new and useful N-phenyloxazolidinone compounds and their preparations, and more particularly to N-phenyloxazolidinone compounds in which the phenyloxazolidinone moiety is linked to a variety of saturated, or partially saturated, 4-8 membered heterocycles containing oxygen, nitrogen, and sulfur through a carbon-carbon bond.\nThe compounds are useful antimicrobial agents, effective against a number of human and veterinary pathogens, including gram-positive aerobic bacteria such as multiply-resistant staphylococci and streptococci, as well as anaerobic organisms such as bacteroides and clostridia species, and acid-fast organisms such as Mycobacterium tuberculosis and Mycobacterium avium. The compounds are particularly useful because they are effective against the latter organisms which are known to be responsible for infection in persons with AIDS."}
{"text": "1. Field of the Invention\nThis invention relates to medical devices, and more particularly, to endoluminal devices and methods for making and using such endoluminal devices.\n2. Description of Related Art\nThe functional vessels of human and animal bodies, such as blood vessels and ducts, occasionally weaken or even rupture. For example, an aortic wall can weaken, resulting in an aneurysm. Upon further exposure to hemodynamic forces, such an aneurysm can rupture, resulting in internal bleeding, and often death.\nVarious interventions have been provided for weakened, aneurysmal, dissected or ruptured vessels, including surgical interventions and endovascular interventions. Endovascular interventions generally include inserting an endoluminal device or prosthesis such as a stent or stent graft into the damaged or diseased body lumen to provide support for the lumen, and to exclude damaged portions thereof.\nThe endovascular prosthesis is delivered in a radially compressed configuration using a catheter delivery system. The catheter is introduced into the lumen system remotely of the repair site and the prosthesis is delivered to the repair site intraluminally. The prosthesis is then expanded to engage the luminal wall. The prosthesis may provide some or all of the functionality of the original, healthy vessel and may further preserve any remaining vascular integrity.\nAn example of a prosthesis that may be used for treating damaged or diseased body lumens is disclosed in PCT Application WO 98/53761, which is herein incorporated by reference. The prosthesis may include a bifurcated stent graft. The stent graft includes a biocompatible graft material and a plurality of longitudinally disposed stents. The stent graft is designed to span and exclude an aortic aneurysm extending between the iliac and renal arteries. Other prostheses that may be used include non-bifurcated stent grafts for spanning and excluding aortic aneurysms within the abdominal aorta or the thoracic aorta.\nOften times, a body lumen may be damaged in an area that includes a branch vessel. For example, there are at least three branch vessels extending from the abdominal aorta, each leading to various body organs. These branch vessels include the celiac, mesenteric, and renal arteries. When an aneurysm includes or is adjacent to one or more of these branch vessels, the prosthesis system must be able to exclude the aneurysm while maintaining fluid flow through the branch body lumen.\nVarious stent grafts have been provided for repairing main body lumens and spanning branch vessels without occluding fluid flow thereto. For example, a main body stent graft may be provided that has one or more fenestrations or apertures in the side wall of the stent graft. The stent graft can be deployed so that the fenestration is aligned with a branch vessel.\nIn many cases, particularly where the damaged portion is positioned at the junction between the main body lumen and the branch body lumen, or where the ostium of the branch vessel is damaged, a main stent graft is insufficient to adequately repair the luminal system. In these situations, it may be preferable to provide a branch lumen prosthesis for positioning within the branch vessel. The branch lumen prosthesis may be used independently, or in conjunction with a main body prosthesis.\nU.S. Published Patent Application Nos. 2005/0222668, 2005/0171598, and 2005/0149166 disclose various systems for repairing branched body lumen systems. Various aspects of each of these disclosures may be used in conjunction with the present invention. U.S. Published Patent Application Nos. 2005/0222668, 2005/0171598, and 2005/0149166 are herein incorporated by reference.\nA branch vessel prosthesis should be capable of complying with a variety of challenging and often competing demands. For example, the branch vessel prosthesis should preferably be highly flexible and capable of tracking through and conforming with a highly tortuous luminal environment. If the prosthesis includes a balloon-expandable stent, the stent should be sufficiently resilient so as not to hinder balloon expansion and/or molding.\nOn the other hand, once the prosthesis is implanted in the body lumen, it must be sufficiently strong and robust to survive a highly dynamic and pulsatile luminal environment that can promote prosthesis damage. This is of particular concern where the branch vessel prosthesis is deployed within a fenestration of a main body prosthesis. During the cardiac cycle, the main body prosthesis will pulse and move with the main body vessel, placing stress on the branch vessel prosthesis at the fenestration. When the main body prosthesis moves, it can exert significant concentrated and localized stresses on the branch vessel prosthesis 11 through the fenestration. Over time, this cyclic wear can cause the branch vessel prosthesis to weaken and eventually to crush under the force of the main body prosthesis, requiring further medical intervention."}
{"text": "1. Field of the Invention\nThe present invention relates to an automatic blood pressure measuring apparatus which can accurately measure a blood pressure of a female subject even at the time of childbirth.\n2. Related Art Statement\nThere is known an automatic blood pressure (BP) measuring apparatus which monitors BP values of a living subject by periodically changing an air pressure of an inflatable cuff applied to a body portion of the subject at a predetermined measurement period and detecting a pressure pulse wave produced in the cuff during the change of air pressure of the cuff. It is known that the BP values measured using the inflatable cuff are reliable.\nMeanwhile, it is medically important to monitor BP values of a female subject at the time of childbirth. However, in the case where the above-mentioned BP measuring apparatus is used to monitor, using the cuff, BP values of a female subject who heavily moves her body each time she feels a labor pain (i.e., uterine contraction), the BP measuring apparatus may not accurately measure the BP values of the subject because of her heavy physical motion.\nIt is therefore an object of the present invention to provide an automatic blood pressure measuring apparatus which can measure an accurate blood pressure of a female subject, even during childbirth.\nThe above object has been achieved by the present invention. According to a first feature of the present invention, there is provided an apparatus for automatically measuring a blood pressure of a female subject, comprising a blood-pressure measuring device which includes an inflatable cuff adapted to apply a pressing force to a body portion of the subject and automatically measures, by changing the pressing force of the cuff, a blood pressure of the female subject at a predetermined measurement period; a labor-pain-signal detecting device which detects, from the female subject, a labor-pain signal which changes in relation to a labor pain of the subject; and a first blood-pressure-measurement starting means for operating the blood-pressure measuring device to start a blood-pressure measurement, when the predetermined measurement period has elapsed and when the labor-pain signal detected by the labor-pain-signal detecting device falls within a reference range.\nIn the present automatic BP measuring apparatus, the first BP-measurement starting means operates the BP measuring device to start a BP measurement, only when the labor-pain signal detected by the labor-pain-signal detecting device falls within a reference range. If the labor-pain signal falls within the reference range, then the female subject is not feeling a labor pain and accordingly she is not moving her body. Since the present apparatus starts a BP measurement when the female subject is not moving her body, it can accurately measure one or more BP values of the subject.\nAccording to a second feature of the present invention, the blood-pressure measuring apparatus further comprises a labor-pain-period-relating-information determining means for determining, based on the labor-pain signal detected by the labor-pain-signal detecting device, labor-pain-period-relating information relating to a labor-pain period at which the female subject periodically feels a labor pain; and a measurement-period shortening means for shortening the predetermined measurement period, when the determined labor-pain-period-relating information indicates that the labor-pain period of the subject has shortened.\nDuring childbirth, it is medically needed to detect more quickly any abnormal tendency of the female subject when the labor-pain-period-relating information indicates that the labor-pain period of the subject has shortened. To this end, the present automatic BP measuring apparatus employs the measurement-period shortening means which shortens, when the labor-pain-period-relating information indicates that the labor-pain period of the subject has shortened, the predetermined measurement period, so that the BP measuring device measures a BP value of the subject at the shortened measurement period. Thus, the present apparatus can more quickly detect the abnormal tendency of BP values of the female subject.\nAccording to a third feature of the present invention, the blood-pressure measuring apparatus further comprises an abnormality judging means for judging, when the labor-pain signal detected by the labor-pain-signal detecting device does not fall within the reference range when the blood-pressure measuring device is performing a first blood-pressure measurement, that the first blood-pressure measurement is abnormal; and a second blood-pressure-measurement starting means for operating the blood-pressure measuring device to start a second blood-pressure measurement when the labor-pain signal detected by the labor-pain-signal detecting device has changed to fall within the reference range again after the abnormality judging means has judged that the first blood-pressure measurement is abnormal.\nWhen the female subject feels a labor pain during a BP measurement of the BP measuring device, the BP measuring device may not accurately measure one or more BP values of the subject. However, the present automatic BP measuring apparatus employs the abnormality judging means which judges, if the labor-pain signal does not fall within the reference range when the BP measuring device is performing a first BP measurement, that the first BP measurement is abnormal, and the second BP-measurement starting means which operates the BP measuring device to start a second BP measurement when the labor-pain signal has changed to fall within the reference range again after the abnormality judging means has judged that the first BP measurement is abnormal. Thus, the present apparatus can accurately measure one or more BP values of the female subject."}
{"text": "A conventional diverter tub spout body 10 is shown in FIG. 1 and comprises a longitudinal portion 22 and a transverse portion 100 which is connected to the top end thereof. A separate wall 23 is connected between the transverse portion 100 and the longitudinal portion 22 and has an connecting tube 24 extending from the separate wall 23, wherein the connecting tube 24 communicates with the interior 31 of the transverse portion 100 so as to be connected to a pipe (not shown). In order to prevent the copper or zinc made diverter tub spout body 10 from getting rust, an electroplating process is necessary which immerses the diverter tub spout body 10 into the electroplating bath solution 11 with the longitudinal portion 22 first to let entire diverter tub spout body 10 immersed below the level of the electroplating bath solution 11. However, referring to FIG. 2, there will be bubbles 12 collected in the recessed portions defined between the connecting tube 24, the separate wall 23 and the spout body of the longitudinal portion 22 so that the electroplating bath solution 11 cannot attach onto the interior of the two recessed portions so that rust will be found in the recessed portions.\nThe present invention intends to provide a diverter tub spout body which has an aperture defined through the separate wall located in either one of the recessed portions so as to dispense bubbles therefrom. The aperture improves the result of the electroplating process so that the shortcomings happened on the conventional diverter tub spout body can be well handled."}
{"text": "Many utility carts exist today which may be contorted into various positions—with each position being particularly adapted for a specific type of use (or to support a load in a particular manner). However, a majority of these utility carts may only be transformed into two or three primary positions. With respect to those utility carts that purportedly may be configurable into more than three positions, such utility carts often do not easily convert from a first position into a second position. Moreover, in many cases, the multi-functional utility carts that are available today are not structured to withstand a substantial load, while simultaneously being capable of collapsing into a compact and lightweight storage position. Still further, the currently-available utility carts are often unstable during use.\nAs explained further below, the present invention addresses many of these shortcomings of the utility carts that are available in the marketplace today."}
{"text": "The present invention relates to novel telomers formed by the reaction of chlorotrifluoroethylene, herinafter designated as \"CTFE\", with fluoroxytrifluoromethane, hereinafter designated as \"FTM\". These telomers have properties which distinguish them from other CTFE telomers and make them commercially attractive products. Such properties include superior solvent characteristics which are useful in formulating non-flammable hydraulic fluids.\nVarious methods of preparing telomers from CTFE are known in the prior art and have been practiced commercially for many years.\nAn article by William T. Miller, Jr. et al in Industrial and Engineering Chemistry, pages 333-337 (1947), entitled \"Low Polymers of Chlorotrifluoroethylene\", describes a process for producing low molecular weight polymers of CTFE by carrying out the polymerization in a solution of chloroform using benzoyl peroxides as polymerization promoters. Other solvents disclosed in the reference as being useful for this purpose include carbon tetrachloride and tetrachloroethylene. The solution is heated in a pressure vessel for 13/4 hours at 100.degree. C., and the unreacted CTFE monomer and chloroform are removed by distillation, leaving a \"crude\" telomer of general formula CHCl.sub.2 (CF.sub.2 CClF).sub.n Cl, which can further heated and distilled to yield products ranging from a light oil to a semisolid wax or grease. In this formula, n is the chain length (the number of repeating units in the telomer chain), and is in the range of 1 to 20.\nA more recent development in this field is described in a series of articles by Y. Pietrasanta et al. entitled \"Telomerization by Redox Catalysis\" appearing in the European Polymer Journal, Vol. 12 (1976). This technology involves the reaction of a chlorinated telogen, such as carbon tetrachloride, with CTFE in the presence of benzoin and a suitable redox catalyst, such as ferric chloride hexahydrate (FeCl.sub.3.6H.sub.2 O). The telomerization reaction is suitably carried out in acetonitrile which is a common solvent for the reactants and catalysts. The telomerization reaction can be illustrated as follows, wherein n is as defined above: ##STR1##\nThe telomers produced according to the above-described process have end groups containing chlorine atoms. When used in a particular application, especially high temperature or corrosive applications, the chlorine atoms can be hydrolyzed, which may result in cleavage of the telomer and a loss of physical properties of the fluid. In order to prevent such a condition, the end groups are stabilized by fluorinating the telomer to replace chlorine atoms with fluorine. The fluorine atoms form a stronger bond with carbon and are less prone to cleavage. This results in a fluid which has superior performance over a wider range of operating conditions.\nFluorination of the telomer is accomplished by reaction with a suitable fluorinating agent, such as chlorine trifluoride or hydrogen fluoride. However, fluorination involves an additional process step which increases processing costs. It would be desirable to reduce the amount of fluorination required or eliminate this procedure in its entirey. It would also be desirable to develop superior telomers having imroved physical characteristics, such as better solvent properties. These are the primary objectives of the present invention."}
{"text": "Due to the large number of components in state of the art electronic circuits, most of their design and production is computer implemented. An important operation performed during design of electronic circuits is timing analysis that validates timing performance of an electronic design. For example, static timing analysis validates timing performance of a circuit design by checking possible paths for timing violations. Performing timing analysis for large electronic circuits can take significant time.\nFurthermore, timing analysis is repeated multiple times for the same electronic design for various combinations of modes and corners. Semiconductor device parameters can vary with conditions such as fabrication process, operating temperature, and power supply voltage. A circuit fabricated using these processes may run slower or faster than specified due to variations in operating conditions or may even fail to function. Therefore timing analysis is performed for various operating conditions to make sure that the circuit performs as specified under these conditions. Such operating conditions for a circuit are modeled using corners that typically comprise a set of libraries characterized for process, voltage, and temperature variations.\nThe timing analysis of a circuit is also repeated for different operating modes, for example, normal operating mode, test mode, scan mode, reset mode and so on. For example, a circuit used in a computer operates in a stand-by mode when the computer is in a stand-by mode. Similarly, during testing phase, a circuit may be operated in a test mode. A mode is modeled using a unique set of clocks, input voltages, and timing constraints in similar operating conditions.\nFor performing timing analysis during implementation and sign-off of a circuit, designers may have to verify a large number of modes and corners. Each circuit design may have tens of modes and tens of corners or more. Since each mode should be verified for each corner condition, the total number of scenarios for which the design is verified is the product of the number of modes and number of corners. This results in the timing analysis being performed a large number of times resulting in high costs.\nOne way to handle the large number of scenarios resulting from multiple modes and corners is to merge the modes into a smaller set, for example, a single mode. Since timing analysis is performed for the combination of modes and corners, reduction in the number of modes reduces the total combinations of modes and corners by a much larger number. For example, if there are 10 modes and 10 corners, the total number of scenarios is 10×10=100. However if the 10 modes were combined to a single mode, the total number of scenarios is reduced to 1×10=10 which is a 90% reduction in the number of scenarios to be verified.\nConventionally, modes are merged manually by designers. However, due to the complexity of constraints associated with circuit designs, the generation of merged modes is difficult to handle manually. Since the number of timing constraints for a given netlist can be large, manual merging of modes can be error prone and have prohibitive costs. For example, a circuit could have millions of constraints and manually determining which modes can be merged together and then merging them may not be practically feasible."}
{"text": "Commercial fishermen catch and process tons of fish a day. Many fishermen have complex, highly-automated processing equipment which dress the fish and produce numerous different commercial products for commercial retailers, such as restaurant chains or grocery stores. Such equipment is highly specialized to perform a specific task, such as cutting the heads and tails off, filleting, deboning, and skiing, etc. One of the reasons such equipment must be tailored to specific tasks is the large variety of shapes and sizes of fish processed, including pollock, salmon, trout, sole, cod, etc.\nTo reduce overall processing costs and to improve quality, commercial fishing vessels include processing equipment and freezers located directly on the ships. As the fish are caught, they are immediately processed and frozen, substantially improving the quality of the final product. This processing equipment also allows the vessels to remain at sea for long periods of time without transporting the fish to a shore-based processing facility. Most of the processing equipment is large, complex, and expensive. It is difficult and expensive to remove the processing equipment from the vessel and replace it with new processing equipment.\nBecause of the expense associated with exchanging the processing equipment, including vessel downtime, most commercial processing ships are intended to catch and process specific types of fish. As an example, a vessel may be fitted with processing equipment designed to process generally oval fish having a relatively thick cross-section, such as pollock or salmon, or generally flat fish having a relatively thin cross-section, such as sole.\nDue to the natural migration of fish, and various fishing laws, individual types of fish are caught and processed during limited parts of the year. Therefore, vessels having equipment capable of processing only one type of fish sit idle during part of the year. If the processing equipment could be readily exchanged or adapted to process both oval fish, such as pollock or salmon, and flat fish, such as sole, the vessels could be operated over a larger part of the year. This would in turn reduce the amount of downtime and associated expense for the vessels.\nOne manufacturer of fish processing equipment is Baader GmbH & Co. KG, Lubeck, Germany (\"Baader\"). One of the fish processing machines produced by Baader is the Baader 182. The Baader 182 is designed to process pollock and salmon shaped fish and includes a conveyor having a plurality of fish holders mounted along the length of the conveyor. Each fish holder includes a recess configured to receive a pollock. While operating, a worker removes fish from a fish hopper located next to the conveyor and places an individual pollock in the recess of each fish holder. The conveyor subsequently moves the fish holders and pollock along the length of the conveyor into contact with rotating blades that cut both the head and tail of the pollock off. The body of the pollock is then carried to additional equipment that cuts the pollock into fillets and removes the bones and skin.\nThe Baader 182 works well with fish having an oval cross-section, such as pollock, but does not allow other types of fish to be processed. The Baader 182 is particularly unsuitable for processing fish with a relative flat cross-section, such as sole. The fish holders used on the Baader 182 are not capable of holding such fish during processing.\nIn order to adapt the Baader 182 to process fish having a relatively thin cross-section, the inventor previously applied for and received a patent on a device for holding flat fish during processing. U.S. Pat. No. 5,358,441 entitled \"Device for Holding Flat Fish During Processing\" issued to the present inventor on Oct. 25, 1994. The device for holding flat fish during processing allows the Baader 182 to process flat fish in addition to pollock. One embodiment of the device included a generally rectangular fish holder that replaced the fish holders on the Baader 182. The fish holders include two arms which extend over a portion of the fins and body of a flat fish and hold it in place on the fish holder. The arms include two recesses that have slanted rear walls that center the flat fish as it is placed within the fish holder. Two slots extend across the width of the fish holder to allow rotating blades to extend beneath the upper surface of the fish holder to ensure that the blades cut cleanly through the thickness of the fish.\nOne of the products that the device for holding flat fish during processing allowed the Baader 182 to be used for is the processing of \"kirimi.\" In order to make kirimi, the head and tall of a sole are removed, leaving the center portion of the sole. The center portion of the sole is then sold to consumers who fry, bake, or otherwise prepare the center portion for consumption. In the past, the center portion was sold complete with any internal organs remaining within the center portion. One of the internal organs remaining in the center portion after removing the head and tall is the kidney, commonly referred to as the \"blood spot.\" The blood spot is located along the backbone of the sole in the internal organ cavity. In the past, the blood spot has been either removed by the individual consumer during preparation or removed by hand during processing after the head and tail are removed.\nIn today's highly competitive fishing industry, it is generally not economical to remove the blood spot by hand, thus the blood spot is generally left in the commercial kirimi product. If the blood spot could be removed during processing, a higher quality kirimi product would result.\nAs can be seen from the above discussion, there exists a need for a method to remove the blood spot and any other internal organs remaining in the kirimi after processing. The present invention is directed toward fulfilling this need."}
{"text": "1. Field of the Invention\nThe present invention relates to a method of fabricating a semiconductor device, and in particular to a method of fabricating a semiconductor device with extremely shallow-junctioned source/drain structures.\n2. Description of the Related Art\nAs the integration level of semiconductor devices increases, the junction depth of the source/drain structures of MOS transistor is reduced more and more, and extremely shallow source/drain junctions are becoming more important.\nReferring to FIG. 1, a conventional method of fabricating a P-type semiconductor device with shallow junctions will be described. In FIG. 1, an N-type well 22 is formed in the semiconductor silicon substrate 21, and then a device isolation layer 23 is formed in the field area of the silicon substrate 21 by LOCOS process. A gate oxide layer 24, gate electrode 25, and insulator spacer 26 are then formed on the silicon substrate 21 and BF.sub.2 ion implantation is performed onto the active area of silicon substrate 21 defined by the field area at the dose of 2E15/cm.sup.2 to 5E15/cm.sup.2 and the energy of 15 keV to 40 keV in order to form P.sup.+ source/drain 27 of the transistor.\nThereafter, an interlayer insulator film 29, e.g. tetraethylorthosilicate glass (TEOS) layer is deposited on the entire surface of the silicon substrate 21 and then a planarizing insulator film 29 is deposited thereon.\nHowever, when the BF.sub.2 ion implantation is performed in order to create the P.sup.+ source/drain 27, the surface layer of the silicon substrate 21 is damaged due to the implanted ion impact so that it becomes an amorphous silicon layer, and some point defects occur under the interface between the amorphous layer and the crystalline layer. Further, during the subsequent thermal process, the diffusivity of boron is increased due to the point defects and therefore the junction depth of the source/drain 27 is increased. In addition, some secondary defects or extended defects occur, and since the electrical activity of the boron is thus decreased, the leakage current in the junction is significantly increased. Therefore, the conventional technique is unsuitable to form an extremely shallow junction with excellent electrical characteristics."}
{"text": "In general, radiation shields which cover portions of the body such as the thyroid area, male gonadal area, female gonadal area, breast area, hands, and eyes are known in the art. Examples of such radiation shields are described in U.S. Pat. No. 4,938,233, filed by Orrison, Jr., issued Jul. 3, 1990, entitled \"Radiation Shield.\" These radiation shields are generally comprised of a body of radiation attenuating material and an attachment member. The attachment member of the radiation shield is configured to be worn around an appendage such as the head, neck, or waist so that the radiation attenuating material protects a particular body portion.\nRadiation shields are intended for use in non-sterile environments during radiological diagnostic procedures or oncological treatments. The shields are worn by patients, clinicians, and physicians to selectively isolate, shroud and protect particular regions of the anatomy. These radiation shields are reusable. The same radiation shield may be worn hundreds of times by different practitioners and patients.\nPrior radiation shields have been permanently coated with vinyl so that the shields may be cleaned after each use. The vinyl coating of the radiation shield is uncomfortable to wear adjacent the skin. The smooth surface of the coating also allows the shield to be easily dislodged once it is placed on the body.\nDespite cleaning, the radiation shields often become soiled with bodily fluids or otherwise unsightly discolorations. These stained radiation shields, which often contain environmentally unfriendly materials, must be thrown away for eventual accumulation in landfills. Further, if the radiation shield is not properly cleaned, the patient or practitioner unknowingly wears a soiled or otherwise unsanitary shield. This is a particular problem for radiation shields which cover areas of the anatomy such as the gonadal region.\nThere is a need for a method which allows radiation shields to be hygienically reused. There is also a need for a low cost reusable radiation shield for protecting particular areas of the body."}
{"text": "Garments, such as pants, skirts, and shirt sleeves, are generally adjusted by stitching or cutting. Adjusting hems on pant legs, for example, requires manual alteration of the garment by cutting the stitches, unfolding the cuff, resetting the cuff to a desired length, and re-stitching the hem. Traditional manual stitching operations are time consuming and prone to errors. Moreover, alterations of garments as described above are permanent unless the garment is altered again.\nUnfortunately, such permanent alterations do not address the many different contexts in which a garment may be worn. For instance, trousers or slacks may be worn with different types of shoes (e.g., high heels, flats, boots, sneakers, etc.). The particular shoe to be worn will, in turn, dictate the desired length of the trousers or slacks. For example, the length of a pair of trousers or slacks will need to be longer for a shoe with high heels than for a flat shoe with little or no heel. A user may have to have the trousers or slacks hemmed by a tailor to the desired length dictated by the particular type of shoe the user intends to wear with the garment. That length, however, may not correspond to the length dictated by a different type of shoe.\nThere have been previous attempts at facilitating garment hem adjustments, which generally included systems embedded in the garment. Some of these attempts were not aesthetically pleasing and/or were difficult to use, while others were uncomfortable for the person wearing the garment.\nIt is, therefore, desirable to have a system for readily adjusting the length of a garment that is efficient, aesthetically pleasing, and comfortable to wear."}
{"text": "In a silver halide photographic element, a color image is formed when the material is exposed to light and then developed using a primary amine color developer. The development results in imagewise reduction of silver halide and the corresponding production of oxidized developer. The oxidized developer then reacts in an imagewise fashion with one or more incorporated dye-forming couplers to form a dye image.\nMagenta dye-forming couplers are employed in subtractive color forming processes. One of the principal coupler types useful for this purpose are those based on a pyrazolone ring. Pyrazolone-based couplers having a coupling-off group linked to the pyrazolone ring by oxygen have long been considered as potentially attractive two equivalent magenta couplers. An oxygen-linked coupling-off group could impart increased activity to the pyrazolone coupler; however the general instability of these couplers toward ambient oxygen makes them difficult to synthesize and impractical for use in a film environment since they decompose during keeping. In particular, pyrazolone couplers having an anilino or acylamino substituent at the 3-position have exhibited unacceptable stability when an aryloxy is employed as a coupling-off group. As a result, pyrazolone couplers have employed either so-called \"four equivalent\" couplers containing hydrogen at the coupling-off position or have employed so-called \"two-equivalent\" couplers containing a coupling-off group having a sulfur or nitrogen atom linked to the pyrazolone ring.\nU.S. Pat. No. 3,419,391 discloses certain types of pyrazolone-based compounds as two-equivalent couplers having high dye-forming reactivity and reduced tendency to form color fog. According to the patent, the pyrazolone ring is not limited to the presence of any particular substituents at the 3-position or elsewhere. Specifically identified substituents at the 3-position include anilino, acylamino, alkyl, amino, alkoxy, amido, carbamoyl, ureido, thio, guanidino, etc. The couplers of the patent may contain an aryl group at the 1-position and, among other things, an alkyl or carboxy ester group at the 3-position. The aryloxy couplers of the patent are said not to produce color fog (printout) and to provide improved reactivity. Thus, they arc said to provide low printout or yellowing in Dmin areas when they are exposed to light or high temperatures, respectively, subsequent to development. No mention is made of the poor keeping of pyrazolone couplers having aryloxy coupling-off groups although their instability is well known in the art.\nU.S. Pat. No. 5,576,167 describes a family of two-equivalent pyrazolone couplers with good stability towards oxygen. The features that make these compounds stable are: (a) an electron-withdrawing CXYZ substituent in the 3-position of the pyrazolone ring, and (b) an aryloxy group, further substituted by electron-withdrawing group(s) in the 4-position of the ring. Although stable, the couplers of this patent require a multistep procedure for their preparation and some of these steps are difficult to perform. None of these references discloses magenta couplers of this invention.\nIt would therefore be desirable to have a pyrazolone-based coupler with an aryloxy coupling-off group present in a photographic element that is stable during synthesis, film manufacture, and during film keeping and would provide an acceptable hue and reactivity but which would not be difficult to prepare. The 1-aryl-3-arylpyrazolo-5-one magenta dye forming couplers of this invention provide these and other desirable advantages."}
{"text": "Field\nApparatuses and methods consistent with the exemplary embodiments relate to a display apparatus and a method of displaying a method thereof. More particularly, the exemplary embodiments refer to a display apparatus connectable to home appliances and a message display method thereof.\nDescription of the Related Art\nHome networks, which have become wide spread in recent years, refer to a home system which enables household electronic and electric appliances to be connected via a single wireless or wire-based system in order to implement two-way communication. The home network allows automated adjustment of not only simple functions of remotely controlling home appliances, i.e., (home automation) but also high-tech functions, including information exchange, monitoring and security through data transmission and reception between home appliances. A user may switch on and off gas or electronic products from outside the home via a computer, a personal digital assistance (PDA) or a mobile phone and access multimedia content, such as education programs, movies, and games, via a TV.\nThe home network may be classified into a data network of home information and communication devices based on a computer, an audio/video (A/V) network of a TV and an audio system, and an information appliances network of household equipment based on white goods, depending on use. In the home network, wire-based technology, such as Ethernet, PLC, IEEE 1394 and home PNA, radio technology, such as IEEE 802.11 WLAN, IEEE 802.15 WPAN and UWB, and home network control middleware, such as UPnP, HAVI, JINI and Home Network Control Protocol (HNCP) are employed, and various kinds of services based on these technologies are being vigorously developed."}
{"text": "1. Field of the Invention\nThe present invention relates to the general art of firearms, and to the particular field of rests and supports for firearms.\n2. Discussion of the Related Art\nIt is well known that hunting is one of the most popular sports in America. Millions of hunters participate in the sport during all seasons. For this reason, the arts associated with hunting and hunting equipment contain many examples of devices and methods intended to make hunting easier and more efficient.\nOne form of hunting requires a hunter to remain stationary for long periods of time. During this waiting time, the hunter is required to be ready to execute a shot whenever a target presents itself. Thus, the hunter must remain stationary for long periods, but ready to shoot at a moment's notice. This may require the hunter to hold his or her firearm in a shooting position for great lengths of time. Not only is this uncomfortable, it may adversely affect the hunter's aim due to fatigue.\nTherefore, there is a need for a support for a firearm that will permit a hunter to remain stationary with the firearm in a firing position.\nWhile the art contains many examples of firearm supports, these supports all have various shortcomings that prevent the hunter from being able to remain stationary for long periods of time yet to have the firearm in firing position ready for an accurate shot. For example, some firearm supports include a tube that is based on the ground. These supports may require the hunter to adopt an uncomfortable position or require the hunter to pick up the firearm to shoot. In the first case, the hunter may not be in the most desirable shooting position when a target appears and in the second case the hunter may lose the opportunity for an accurate shot during the time the firearm is being removed from the support and/or moved into shooting position. Other firearm supports attach the firearm to a support, such as a tree or even to the hunter's lap. These supports do not permit accurate and rapid shooting.\nTherefore, there is a need for a support for a firearm that will permit a hunter to remain stationary with the firearm in a firing position while holding the firearm in a shooting position during the wait.\nAiming a firearm is never easy, and accurately and quickly aiming a firearm is often what separates a successful shot from an unsuccessful shot. As discussed above, while the art contains many examples of firearm supports, these known firearm supports are not amenable to accurate and rapid aiming of a firearm after a long period of firearm support.\nTherefore, there is a need for a support for a firearm that will permit a hunter to remain stationary with the firearm in a firing position while holding the firearm in a shooting position during the wait and will permit the hunter to quickly and accurately aim the firearm for a shot."}
{"text": "1. Field of the Invention:\nThe present invention relates to a control system using a plurality of computers and, more particulary, to a control system using a plurality of microcomputers for controlling a copying machine main body and its peripherals.\n2. Description of the Prior Art\nRecently, tremendous advances have been made in electronic techniques. In particular, with the recent trend toward higher integration of electronic circuits, high-performance one-chip microcomputers are used for a variety of equipment. Microcomputer-aided control of a copying machine is becoming popular. With requirements for higher performance, higher operation speed and multipurpose operability, a plurality of microcomputers are used in place of a single microcomputer so as to control the copying machine main body and its peripherals (sorter, ADF and the like) separately in accordance with different functions. For this purpose, parallel or serial data transfer through intelligent terminals using I/O ports of the microcomputers is being developed.\nIn a system wherein control power supplies for the respective microcomputers are separately provided, problems due to the rise of the control power supply voltage at the initial power ON, that is, problems of operation timing of the respective microcomputers, or problems of data transfer error due to the states of the I/O ports during the initial reset (initialization) of the respective microcomputers, are encountered during data transfer. The copying machine main body or any one of its peripherals may as a consequence operate erratically."}
{"text": "1. Field of the Invention\nThis invention relates to a serial type recording apparatus and method which is applicable to a recording apparatus in a business machine or a communication instrument such as a copying machine, a facsimile apparatus, a personal computer or an office computer, and in particular to a recording apparatus and method which can be effectively applied to an ink jet recording apparatus in which a carrier is caused to scan in the direction of a recording column to thereby accomplish recording.\n2. Related to Background Art\nAmong the serial type recording apparatuses, such as ink jet recording apparatuses, wire dot printers and heat transfer printers, proposed before the present invention, there is one in which a recording head is carried on a single carriage and when image recording using different ink is to be effected, the recording head is interchanged with another recording head by a manual operation. However, in such an apparatus, it is not intended that a plurality of recording heads be carried on the carriage to effect color recording, and such apparatus is not constructed so that a plurality of recording heads can be carried on the carriage. In any case, black documents and the black printing mode data processing are more often used for printing in a terminal instrument than colored graphics and colored documents.\nInstead of the manual operation, it is conceivable to move a carriage carrying a plurality of recording heads thereon serially to thereby enable printing in any color to be accomplished, but this will result in an unused construction in which a recording heads are always moved or continue to be moved in a state capable of recording. Accordingly, in the heretofore known ink jet recording, a recording heads are subjected to the recovery process or ink is heated, which leads particularly to the waste of time and the waste of electric power. Also, the carriage carries the maximum weight thereon at all times and, therefore, the load for moving the carriage becomes great.\nOn the other hand, a color ink jet printer carried recording heads for four colors, i.e., Bk (black), Y (yellow), M (magenta) and C (cyan) in parallel on a single carriage and has effected printing by causing the carrier to scan.\nHowever, in the ink jet printer, during black printing, the other printing heads are exposed to the atmosphere, and this causes problems such as clogging and adherence resulting from the desiccation of ink. These problems arise earlier in a recording head not being used than in a recording head being used, and even if the recovery process is carried out, the period during which these problems arise becomes very short. During recovery process, it is usually the practice to return the carriage periodically to its home position even during black printing to thereby effect idle discharge, or to effect suction of ink from the color printing heads even during the closing of the main switch or the suction recovery operation,. and this brings about the inconvenience that the recording time is prolonged. Also, by this recovery process, the color inks are wastefully consumed, and this has been uneconomical.\nAlso, even when high-speed printing is to be effected, it has been necessary to drive the heavy carrier carrying the four recording heads thereon and thus, an expensive and bulky motor for driving the carrier has been necessary.\nEven in printing conditions under which black printing can be executed, the absence of other color inks makes printing impossible.\nAlso, in spite of the execution of printing being possible if the temperature control of only the black head is effected in printing conditions under which black printing can be executed, waste such as affecting the temperature control of the other color heads, for example, warming all the heads so as to assume 35.degree. C. or higher, is unavoidable.\nThe present invention has been made on the basis of the above-described background art, and more particularly from a point of view which has not heretofore been foreseen."}
{"text": "This application incorporates by reference Taiwan application Serial No. 090125720, filed Oct. 17, 2001.\n1. Field of the Invention\nThis invention relates in general to a controlling device for a motor. More particularly, the invention relates to a controlling device for a motor capable of being driven by an AC power source or a DC power source.\n2. Description of Related Art\nA conventional AC (alternative current) fan uses an AC power source as its working power source, so that the coil wiring of the stator or the rotor is to form an induced magnetic field to operate the motor. However, since the rotational speed of the motor is proportional to the frequency of the AC power source, the rotational speed of the motor cannot be increased due to too low frequency of the AC power source. For some situations requiring a high wind flow, the motor cannot increase its rotational speed because the conventional fan\"\"s motor uses the AC power source. The formula of fan speed expressed as follow: r.p.m.=(N*f/2)*(60 sec.), wherein r.p.m. is the round-per-minute speed, N is the pole of wirings, and f is the frequency of AC power source.\nConventionally, in order to obtain a high wind flow by using the AC motor, multiple motor sets are combined to achieve an output of high wind flow. An alternative method is to increase the number of pole wirings within the AC motor, so that the motor\"\"s rotational speed is increased to arise the wind flow. However, to increase the number of pole wirings is not only that the coil wirings are very difficult, but the motor\"\"s manufacturing cost is also increased. Therefore, these solutions do not meet the economic requirements.\nIn addition to the above defects, the inductance of coil is always in an high impedance if the fan working in AC mode, which is because the current of coil is alternatively changed. The working current of the AC motor can be expressed as followed.\nI=V/Zxe2x80x83xe2x80x83(1)\nWhere I is the working current of the AC motor, V is the voltage provided by the power source, and Z is the impedance of the coil. From the formula (1), the impedance Z of the coil is expressed below.\nZ={square root over (R2+XL2)}xe2x80x83xe2x80x83(2)\nwhere XL is an AC inductance, XL=2xcfx80f L. Therefore, when the motor coil is from the static status to the operation status, it needs a sufficient current to make the coil to create a sufficient starting torque. In addition, the current is proportional to the voltage, and therefore, the starting voltage for the motor is also very high. Accordingly, the low voltage characteristic of the AC motor is very worse. Moreover, when the conventional AC motor is locked, because no break protection device, the temperature increases will damage the fan.\nAs a result, for the current fan using the AC motor, its rotational speed, starting up approach, protection device and manufacturing cost are restricted by the AC power source, and therefore cannot overcome the aforementioned defects. Therefore, to further improve the AC motor structure is demanded.\nIn light of the foregoing description, it is an object of the present invention is to provide a motor controlling device, wherein an AC power source is used and the rotational speed of the motor of the fan can be effectively increased.\nAccording to the object(s) mentioned above, the present invention provides a motor controlling device for an AC/DC driven motor. This controlling device is used for an AC/DC driven motor. The motor comprises plural coils, and the motor controlling device comprises an AC/DC converting unit, a voltage stabilizing unit, a phase detecting/controlling unit, and a high voltage driving unit. The AC/DC converting unit is inputted an AC power source to convert the AC power source into a high voltage DC voltage. The voltage stabilizing unit is coupled to the AC/DC converting unit, for receiving the high voltage DC voltage to convert it into a low voltage DC voltage. The phase detecting/controlling unit is coupled to the motor and the voltage stabilizing unit, for receiving the low voltage DC voltage as a work voltage, and for detecting a polarity of the permanent magnet to output a first detecting and a second detecting signals accordingly, so as to provide a shut down protection when the motor is locked. The high voltage driving unit is coupled to the AC/DC converting unit, the phase detecting/controlling unit, and the motor, for receiving the high voltage DC voltage as a work voltage, and for receiving the first detecting and the second detecting signals to output accordingly a first driving and a second driving signals so as to drive the motor."}
{"text": "1. Technical Field\nThe present invention relates to medication information and alarm devices that are programmed with medication information and timing sequences for the proper administration of the medication contained within.\n2. Description of Prior Art\nPrior art devices of this type have relied on a wide variety of different medication systems that provide audible and visual alarms to the patients to indicate proper dosage times as well as last access time to the medication container. This type of information and indicators are needed considering the number of patients that are infirmed and take multiple drugs which may inter-react with one another if not taken in accordance with prescription instructions. While medication bottles have labels stating the patient's name, type of drug, dosage and any associated warnings, many patients are still unaware of the contents of their prescription container. This is especially true for the elderly and infirmed or moderately mentally disabled patients. It is therefore a major problem with these individuals to comply with their treatment directions. Some patients cannot readily understand and act on the label information or alternately they forget to take their medication at the proper time or skip a dosage which may be critical to the effectiveness of the medication.\nPrior art patents have attempted to respond to this problem with a number of medication alarms and visual reminders, see for example U.S. Pat. Nos. 4,448,541, 4,572,403, 5,347,435 and 5,495,961.\nIn U.S. Pat. No. 4,448,541 a medical timer apparatus is disclosed having a magnetic and pressure activated switch that once activated indicates a timed alarm unit on the medication container. This device is activated by the patients gripping the container or opening it, resetting the time to the next dosage alarm required.\nU.S. Pat. No. 4,572,403 is directed to a timed dispensing system in which multiple dosages of medication can initially be placed and that are accessible only at the appropriate time and the appropriate amount.\nU.S. Pat. Nos. 5,347,453 and 5,495,961 disclose and claim portable programmable medication alarms and proper dosage use requirements. These devices use audible alarm signals and generate graphic representation of the prescribed administration time, dosage amount and medical instructions. The devices are programmable by a central computer through an interface device."}
{"text": "Software publishers often attempt to restrict access to portions of compiled software executables to thwart would-be reverse engineering while still allowing the executables to function properly. Reverse engineering is the practice of dissecting and/or analyzing software to understand how it works. On certain systems, reverse engineering can retrieve information stored within software such as information related to cryptographic keys or copy protection schemes. Reverse engineers can even tamper with the software itself.\nWould-be attackers typically have access to the executable code for software and can try to reverse engineer it at their leisure. Those of skill in the art are aware of many approaches to harden a piece of software to such reverse engineering attempts. However, problems with these approaches include rapid growth of source code and increased computational requirements. Such source code can be very difficult to maintain. Further, if an attacker compromises the hardening scheme, the software is difficult to rewrite.\nSuppose a software publisher includes proprietary keys, algorithms, and copy protection related information in software for sale or for free. Reverse engineers can quickly obtain legitimate or illegitimate copies of the software and commence testing the software for weaknesses and thereby extract proprietary keys and algorithms.\nIn modern cryptography, a common methodology is the use of mathematical problems that are believed to be difficult to solve, in order to give bricks to design strong systems. Amongst these problems, one may find the factorization or the discrete logarithm problem. This methodology has been shown to be very efficient, and was one of the key elements to transform the art of cryptography into a real science.\nThe most important success of basing cryptography on hard problems is the arrival of proven security: with some mathematical proof and rigorous analysis. It is possible to prove certain security features of a scheme, supposing only that these problems resist to algorithms in practice. White box cryptography in general lacks sufficient security compared to regular black box cryptography, where huge improvements were made to support such security in the standard model. It is indeed hard to prove something in the white box cryptography case because attackers have access to internal states or values."}
{"text": "1. Field of the Invention\nThe present invention relates to a circuit board assembly including at least one circuit board with at least one embedded optical waveguide which has at least one locally bent part reaching at least an upper/lower external surface of the circuit board to enable its coupling to an optical input/output device.\n2. Background Information\nSuch a circuit board assembly is generally known in the art, e.g. from the U.S. Pat. No. 3,872,236.\nIn this known assembly the waveguide is an optical fiber which is embedded in a substrate, the bent part of which reaches the surface of the broad, this bent part being polished or cut to form a tap enabling its coupling with the input/output device.\nLight transported by the fiber either leaves the latter at the location of the tap under a well defined angle which depends on the curvature of the fibre, or is injected therein under a well-defined angle. As a consequence, the input/output device has to be coupled to the fiber either by abutment contact with the tap, or by using optical devices such as lenses, prisms, holograms etc, to guide the light from/to the input/output device to/from the fiber tap at the required angle."}
{"text": "In a textile processing plant, various types of textile machines are usually arranged in a row. In a spinning mill, these are for example spinning or washing machines. To clean these machines, they have to be blown out at various points. For example, to clean the whorl valve of a spinning machine, a blast of air current must be directed to the whorl valve by means of a nozzle.\nIn order to be able to blow out the various machines at different points, it is known from German Patent No. 34 25 545 to direct two hoses from a distributor station which have different openings to the machines. These two hoses are located on two horizontally extending connecting pieces on the distributor station and are guided vertically downward along a curve of 90.degree.. A flat air-conducting element, which is swivelable about a horizontal axis and seals either the one or the other connection orifice is located in the distributor box.\nThe device has the disadvantage that as one of the two hoses is always acted upon by compressed air, it is not possible to separate the two hoses from the compressed air source. Furthermore, only two compressed air hoses can be connected to the distributor station.\nGerman Patent Publication GM 85 08 228 describes a cleaning device which can travel along the textile machines, having a ventilator on which at least two hoses are located as well as a distributor box for distributing the blast-air to the two hoses. Within the distributor box, there is a valve with a horizontal swivel axis which leads the blast of air to one of several blast hoses attached to the distributor box. One of the blast hoses is used to blow on a washing machine, the other to blow on a spinning machine and a further blast hose is provided for blowing out the whorl valve of a spinning machine.\nA disadvantage of this device is that the distributor box must be constructed quite high due to the arrangement of the valves. Moreover, the number of hoses to be charged with air is limited to three. Finally, there is no position in which none of the hoses is charged with air.\nGerman Patent DE 36 18 934 A1 describes a device for cleaning textile machines having blowers guided above the textile machines, to whose pressure side two vertically extending blast hoses, which have different openings directed to the machine, are attached via a distributor box. In the distributor box, there is an air-conducting element which conveys the blast air to one or the other blast hose. This air-conducting element is a flat baffle which is located inside the distributor box above the connection openings to which the blast hoses are attached. In one position of this baffle, the one connection is covered and, in the other position, the other connection is covered. As a result, the air-blast hose located at the connection which is covered in each case is shut down.\nThis arrangement does enable the simultaneous sealing of all blast-air hoses, however, only if they are not too numerous. With more than three blast-air hoses, the arrangement can no longer be used. Moreover, the air-conducting baffle must be moved electromotively, as it is not manually accessible inside the distributor box."}
{"text": "Unless otherwise indicated herein, the approaches described in this section are not prior art to the claims in this disclosure and are not admitted to be prior art by inclusion in this section.\nThe 3rd Generation Partnership Project (3GPP) is a globally applicable third generation mobile phone system specification that is a result of collaboration between various groups of telecommunications associations, including the European Telecommunications Standards Institute, the Association of Radio Industries and Businesses/Telecommunication Technology Committee (ARIB/TTC), China Communications Standards Association, and the Alliance for Telecommunications Industry Solutions. 3GPP work is ongoing with Universal Terrestrial Radio Access Network (UTRAN) Long Term Evolution (LTE). Long Term Evolution (LTE) is a next generation broadband communication technology developed by the 3rd Generation Partnership Project (3GPP) in order to meet high-speed requirements. A LTE system may realize high-speed packet-based communication at up to 100 Mbps in the downlink.\nIn 3GPP technical specification of Rel-8 or Rel-9, two kinds of radio resource measurement need to be performed by a User Equipment (UE), one is Radio Link Monitoring (RLM) measurement, and the other one is Radio Resource Management (RRM) measurement.\nA heterogeneous environment is under discussion in LTE-Advanced (LTE-A), in which environment a plurality pico cells operating with relatively low transmission power are distributed in a macro cell operating with relatively high transmission power. For a UE located at the coverage boundary between a Macro cell and a Pico cell, if the received downlink signal from the Macro cell is stronger than that from the Pico cell, the UE is associated with the Macro cell and receive downlink signals from the Macro cell. While for the uplink, since the Pico cell is closer to the UE than the Macro cell, its is better select the Pico cell as the reception point of the UE, in order to reduce the transmit power of the UE and the interference to normal UEs in the Pico cell. Therefore, for this UE, the downlink serving cell is the Macro cell but the uplink serving cell is the Pico cell, which is called uplink and downlink decoupling.\nFIG. 1 illustrates a scenario of a typical decoupling handover process. As shown in (A) of FIG. 1, Macro evolved NodeB (eNB) serves UE at both the uplink and downlink at the beginning. The UE moves around in the system. When the UE moves such that Pico eNB becomes closer to the UE than Macro eNB, the UE switches uplink bearers from Macro eNB to Pico eNB first, as shown in (B) of FIG. 1. If both the Macro eNB and the Pico eNB support decoupling, the UE will maintain its downlink carriers to Macro eNB and its uplink carriers to Pico eNB at the area where the signal power from Pico eNB is lower than that from Macro eNB. If the UE moves to be further closer to Pico eNB and that the signal power from Pico eNB is larger than that from Macro eNB, the UE switches its downlink bearers from Macro eNB to Pico eNB, as shown in (C) of FIG. 1.\nThe same criterion is applicable to the cell change from Pico to Macro. The downlink carriers are switched from Macro eNB to Pico eNB first when the signal power from the Macro eNB is larger than that from Pico eNB. After that, the uplink carriers are then switched from Macro eNB to eNB if needed. That is, when a decoupling handover occurs to the UE, it switches one-direction carriers first, and switches the other direction carriers later.\nTo activate or deactive uplink and downlink decoupling, the network has to rely on RRM measurement and report mechanism.\nThe current RRM measurement and report mechanism does not differentiate a decoupling handover case and a normal handover case, i.e., uplink and downlink carriers are switched together. In order to support the decoupling, the UE shall be configured to report the measurement when either the criteria for a normal handover or the criteria for a decoupling handover is satisfied. When a normal handover is to be decided, the measurement reported when the criteria for a decoupling handover is unuseful. Similarly, when a decoupling handover is to be cided, the measurement reported when the criteria for a normal handover is unuseful. Therefore, unnecessary Radio Resource Control (RRC) report is generated, which wastes the bandwidth resource of the network and the processing resource on the UE."}
{"text": "Gate valves are commonly used in the mining, pulp and paper, and food handling industries, among others. As known in the art, the gate moves linearly, perpendicular to the direction of flow. Typically, gate valves are used in a fully opened or fully closed position, and generally considered not suitable for throttling applications. In order to seat properly when fully closed, the typical gate valve is shaped to occupy all or substantially all of the bore. While this provides for robust sealing, the shape of the gate also has its disadvantages.\nIn a fully open position, areas within the bonnet or guide of a typical gate valve, known in the art as the chest area, may be exposed to flowing materials. When used in higher flow applications, this flow can have a detrimental effect on the integrity of the valve. Materials can also get clogged within the chest area of the valve body and the guide. Even if the chest area could be suitably sealed off, the traditional shape of the gate may create flow characteristics not desirable for throttling applications."}
{"text": "Many therapeutic treatments are administered to a patient while they are sleeping or are attempting to fall asleep. While these treatments may achieve their intended result, they also often severely affect the quality of sleep that the patient gets while undergoing these treatments. These treatments often interrupt the patient's normal progression of sleep, causing transient arousals. While these arousals do not result in the awakening of the patient, they often pull patients from deeper stages or higher quality states of sleep. Patients often do not reenter these deeper stages of sleep for a relatively long period of time.\nIn some instances, a therapeutic treatment may cause numerous arousals. This fragments the patient's sleep and prevents the patient from reaching the deeper stages of sleep. Studies have shown that fragmented sleep results in excessive daytime sleepiness. This, in turn, is a direct contributor to many accidents, to a general feeling of lethargy, deterioration of cognitive performance, and/or daytime sleepiness, in the patient.\nOne example of therapeutic treatments causing sleep fragmentation is in the treatment of sleep disorders. Continuous Positive Air Pressure (CPAP) treatments are a primary remedy for a number of sleep disorders such as sleep apnea, hypopnea, and snoring. CPAP treatments consist of delivering a constant positive airway stream of air pressure into a patient's airway during sleep in order to keep the patient's airway from collapsing upon itself. State-of-the-art CPAP machines, often called auto-titration PAP (APAP) machines, automatically adjust the pressure of the delivered air in order to accommodate a patient's respiratory pattern. to the rapid changes of pressure in the patient's airway caused by the APAP machines. Another drawback of current state-of-the-art APAP machines is that they are subject to either false positives (such as when UAR and/or natural irregular breathing events are not pre-empted or do not occur, despite false detection of such and associated treatment control change) or false negatives (such as when genuine upper airway resistance (UAR) and/or related events are pre-empted or do occur but are not detected or responded to with treatment control change). This is due in part to the reliance of these machines on the correct interpretation of an inspiratory waveform and the inaccuracies related to the interpretation of the underlying waveform by the APAP machine. This can also be due to current state of the art gas delivery (or other treatment control such as pacemaker devices) devices inability to enable suitable algorithms to detect and adapt their computation detection sufficiently to pre-empt or predict the probability or onset likelihood of shallow breathing, UAR, arousals, and or associated sleep fragmentation or sleep quality deterioration.\nThe inspiratory waveform varies periodically for reasons not always associated with upper airway resistance. The use of inspiratory waveform as the primary or only means of detection of UAR-related events can cause remedial auto titration measures to be taken when none should be. This is particularly evident where the inspiratory waveform analysis technique does not employ an underlying time-course computational method. The time-course computational method refers to comparing a previous sequence of breaths (prestored from previous treatment session or stored from current session breathing data) or the current breath and comparing the variations or changes as an inferred measure of arousal or sleep fragmentation onset. Excessively rapid or excessively insensitive pressure changes often occur when an auto-CPAP machine tries to correct a normal non-UAR related event, or misses detecting the presence of subtle shallow breathing, hypopnea or UAR, respectively. It is believed that the primary cause of sleep fragmentation is the rapid pressure changes in the patient airway produced by the current APAP machines.\nIn addition to the above, studies have also suggested that some APAP machines are limited in their ability to accurately detect the onset or incidence of shallow breathing, mild hypopnea, or UAR events. This limitation is also possibly attributed to limitations of the machines in interpreting the wave form. Misdiagnosis of such mild hypopnea events results in increased UAR which in turn results in arousal and subsequent sleep fragmentation.\nCurrent state-of-the-art therapeutic devices do not optimally adapt to minimize arousals during therapy. Each patient's arousal threshold is affected by varying parameters, yet current state of the art devices do not have adaptive control algorithms that can adapt their treatment levels to accommodate a number of these varying parameters. These varying parameters include (but are not limited to) sleep history such as sleep deprivation or sleep propensity, physiological factors, psychological factors including (but not limited to) stress or anxiety, environmental factors including temperature; noise; lighting; vibration, factors such as varying threshold to arousals with changing age, drugs and alcohol effects to arousal thresholds and others.\nConsequently, in light of the inherent drawbacks in current therapeutic methods for administering treatments to patients who are sleeping or are attempting to sleep, there exists a need for an apparatus and method of monitoring for patient arousal and for adapting a therapeutic treatment to minimize arousal."}
{"text": "The present invention relates to a process for the production of rumen bypass feed supplements. The process converts glyceride oils to their respective fatty acid calcium salts. In particular, the present invention relates to a process for the production of calcium salts of unsaturated fatty acids derived from fish oil. The calcium salts of the present invention, when fed to cattle, provide reproductive benefits, in particular, an increase in fertility as embodied in an increased rate of impregnation. The present invention therefore also relates to methods for providing such benefits in a ruminant.\nDairy cows must be impregnated once a year to maintain a lactation cycle in which milk is produced for ten months at a time with two month rest periods in between during which the cow is dry. Given the gestation period of a dairy cow, the objective is to impregnate the cow within 83 days after calving. The efficient management of a dairy herd thus requires that the cows be maintained at the peak of fertility to ensure re-impregnation within 83 days.\nAccordingly, there exists a need for nutritional supplements that promote dairy cow fertility. Fish oil fatty acids have become the focus of numerous research programs that seek to capitalize on their nutritional and physiological properties. WO 99/66877 discloses the use of omega-3 fatty acids of fish oil origin to increase fertility in animals including cattle. Among the omega-3 fatty acids disclosed are eicosapentaenoic acid (EPA) and docosahexaenoic acid (DHA).\nUnsaturated fatty acids, however, undergo hydrogenation to saturated fatty acids by microbial action in the rumen and must be fed to ruminants in a protected form. The most familiar form in which fatty acids in general are protected from microbial action in the rumen are the fatty acid calcium salts disclosed by U.S. Pat. Nos. 4,642,317; 4,826,694, 4,853,233; 4,853,233 and 4,909,138. This form of fatty acid protection is widely accepted in the dairy industry.\nFish oils have a glyceride content of 100%. That is, all of the fatty acids in fish oils are in the glyceride form. Fatty acid glycerides do not readily react to form calcium salts using the processes disclosed by the above-listed patents. For a product to be commercially feasible, glyceride levels below about 5 weight percent are desirable to produce a free-flowing and stable product.\nU.S. Pat. No. 5,382,678 discloses a process that reportedly can be used to prepare fatty acid calcium salts from feedstocks having glyceride contents as high as 40 weight percent, with the resulting product having a residual glyceride content of less than about 5 weight percent. Products with residual glyceride contents above 5 weight percent lack storage stability, and are susceptible to oxidation, post-heating, melting, subsequent product solidification, and a tendency to form lumps upon storage. Under industrial conditions, however, it has not been possible to consistently obtain residual glyceride levels below 5 weight percent once the initial glyceride content of the fatty acid feed stock is above about 25 weight percent when using the process of U.S. Pat. No. 5,382,678.\nHydrolyzing the glycerides to levels below 25 weight percent is not commercially feasible. Commercial omega-3 fatty acids in the free fatty acid form are so costly as to be commercially unfeasible. One can reduce the glyceride content of the fish oil starting material by blending it with a low glyceride content fatty acid feedstock, such as Palm Fatty Acid Distillate (PFAD), which has a glyceride content of about 15 to 20 weight percent.\nHowever, the quantity of PFAD that would have to be added to fish oil to reduce the glyceride content to levels commercially feasible for use with the process of U.S. Pat. No. 5,382,678 dilutes the concentration of desirable omega-3 fatty acid such as EPA and DHA to ineffective levels. That is, the levels of DHA and EPA in the resulting calcium salt are so low that quantities of calcium salt must be added to the daily feed ration at levels above what is considered acceptable by the dairy industry.\nTo be commercially viable, omega-3 fatty acid calcium salt feed supplements must have DHA and EPA concentrations high enough to confer the beneficial effects of these omega-3 fatty acids when quantities of the calcium salt are added to feed ration at levels considered acceptable to the cattle industry. Therefore, a need exists for a process by which calcium salts may be prepared from fish oils having high concentrations of omega-3 fatty acid with the calcium salts produced with reduced levels of unreacted glycerides in a free-flowing and stable form easily handled by customers.\nThis need is met by the present invention. It has now been discovered that fatty acid calcium salts having acceptable levels of residual glycerides can be prepared from high glyceride content starting materials by using elevated levels of calcium oxide, making it possible to prepare fatty acid calcium salts from feedstocks containing levels of fish oil effective to provide useful concentrations of omega-3 fatty acids in the finished product.\nTherefore, according to one aspect of the present invention, a method is provided for the preparation of fatty acid calcium salts, which includes the steps of:\n(a) providing a fatty acid feedstock having a glyceride content between about 30 and about 60% by weight;\n(b) adding to the feedstock from about 2 to about 3 equivalents of calcium oxide relative to the feedstock, so that a reactive admixture is formed; and\n(c) adding to the reactive admixture from about 2 to about 5 equivalents of water relative to the calcium oxide, so that the calcium oxide hydrates and neutralizes the fatty acids to form calcium salts.\nThe method of the present invention thus includes the use of feedstocks derived from fish oils diluted to glyceride contents between about 30 and about 60 weight percent with low glyceride content fatty acid feedstocks, such as PFAD. Other suitable sources of low glyceride content fatty acids include fatty acids from soy, cottonseed, corn and other vegetable fatty acid distillates, tallow, yellow grease or other animal or fish derived free fatty acid sources produced by deodorization, refining, hydrolyzation or other processes common in the fats and oil industry.\nThus, the method of the present invention obtains fatty acid calcium salts having useful concentration of omega-3 fatty acids and acceptable levels of residual glycerides that heretofore could not be obtained on a commercial scale using prior art manufacturing techniques. Therefore, according to another aspect of the present invention, fatty acid calcium salts are provided containing omega-3 fatty acids and residual glyceride levels below about 5 weight percent that are prepared by the method of the present invention. More specifically, a fatty acid calcium salt product is provided having a residual glyceride content below about 5 weight percent and containing from about 1 to about 10% by weight of EPA calcium salt and from about 1 to about 10% by weight of DHA calcium salt.\nThe DHA- and EPA-containing fatty acid calcium salts enhance the fertility of ruminants without using 100% glyceride content fish oil feedstocks. That is, beneficial results are obtained from feedstocks blended with fatty acids from sources other than fish oils.\nTherefore, according to still yet another aspect of the present invention, a method is provided for increasing fertility in a ruminant, in which the ruminant is fed an effective amount of the EPA- and DHA-containing fatty acid calcium salts of the present invention. The method of the present invention is particularly effective to enhance the fertility of female ruminants, especially dairy cows. Methods in accordance with the present invention begin feeding the supplements daily to a female ruminant from about 21 days before to about 28 days after parturition and feeding continue at least until conception occurs. The fertility enhancement obtained by the calcium salts of the present invention also includes a reduction in embryonic death in the months following conception. Therefore, methods in accordance with the present invention continue feeding the supplements to a female ruminant for at least 30 days, and preferably for at least 60 days after conception.\nThe above and other features and advantages of the present invention will become clear from the following description of the preferred embodiment.\nThe present invention provides a process by which high glyceride content fatty acid feedstocks may be converted to free-flowing powder or granular fatty acid calcium salt rumen bypass feed supplements, making it possible to prepare fatty acid calcium salts from fish oils, wherein the calcium salts contain useful and beneficial quantities of omega-3 fatty acids. The use of such high glyceride content fatty acid feedstocks represents a significant departure from conventional processes for the manufacture of fatty acid calcium salt feed supplements. The term xe2x80x9cglyceridexe2x80x9d as employed herein includes C10-C22 fatty acid monoglycerides, diglycerides, triglycerides, and any mixture thereof.\nIn a typical process according to the present invention, fatty acid feedstocks are added to a production vessel. The mixing should be accomplished in a kettle designed so that intensive and intimate contacting occurs between the calcium oxide and fat/oil admixture so that a homogeneous dispersion of the calcium oxide particles results. The types of internal mixing elements span a wide gap but would include those with propeller, turbine, plows with chopper blades, or preferably xe2x80x9cCowles-typexe2x80x9d dispersing blades as examples, but others may apply. These same devices would also be suitable for dispersing and homogenizing the water fraction into the fat/oil-calcium oxide admixture.\nFatty acid feedstocks are employed containing up to about 60 weight percent of the fatty acid content in the form of fatty acid glycerides. Glyceride levels between about 30 and about 60 weight percent are preferred.\nSuitable feedstocks include essentially any glyceride or glyceride derivative having a fatty acid profile determined to be nutritionally or physiologically beneficial to a ruminant. Beneficial fatty acid profiles are readily identified by those of ordinary skill in the art, and may be derived from any known source, inclusive of fatty acid sources of animal, vegetable or fish origin. This includes distillates and soap stocks of lard or tallow, vegetable oils such as canola oil, sunflower oil, safflower oil, rapeseed oil, soy bean oil, olive oil, corn oil, palm oil, and the like, and byproducts thereof, as well as fish oils and byproducts thereof.\nSuch fatty acid feedstocks typically contain from about 10 to about 100 weight percent of the fatty acid content in the form of fatty acid glycerides, from about 0 to about 90% by weight of free fatty aids, and less than 5% by weight of moisture, insolubles and unsaponifiables. When necessary, the glyceride content may be reduced to about 60% by weight and lower by adding fatty acid distillates such as PFAD to the feedstock or by pre-treatment to convert a portion of the glycerides to free fatty acids, either catalytically through the use of enzymes, including lipases, or by hydrolysis. Pre-treatment processes may also include processes that increase the level of desirable fatty acids, for example, cold acetone extraction may be used to increase the level of EPA and DHA in fish oil.\nThe present invention thus makes possible the preparation of fatty acid calcium salts from fish oils, which have a 100% glyceride content. According to one embodiment of the present invention, from about 15 to about 50 weight percent of fish oil is blended with from about 85 to about 50 weight percent of a fatty acid feedstock sufficiently low in glyceride content to provide a blend with a glyceride content of less than about 60% by weight. Blends of fish oils with PFAD within these weight ratios are included within the scope of the present invention. A blend containing from about 20 to about 35% by weight of fish oil is preferred.\nThe process of the present invention is particularly well suited for the preparation of fatty acid calcium salts containing beneficial levels of DHA and EPA from fatty acid feedstocks containing high levels of fish oil. Oils containing from about 7 to about 16% by weight DHA and from about 10 to about 17% by weight EPA are preferred.\nThe feedstock may also contain up to about 95% by weight of unsaturated C:16-C:22 fatty acids. Unsaturated fatty acid levels between 50 and about 80 weight percent are preferred. In general, unsaturated fatty acids having from 16 to 22 carbon atoms and from 1 to 6 double bonds are suitable for use with the present invention. Polyunsaturated fatty acids are preferred, with examples of desirable polyunsaturated fatty acids including fish oil-derived omega-3 and omega-6 fatty acids. Suitable fish oil sources include menhaden, herring, mackerel, caplin, tilapia, tuna, sardine, pacific saury, krill, and the like.\nIt may be necessary to heat the fatty acid feedstock to form a uniform, liquid admixture, depending upon the degree of saturation. A temperature up to about 175xc2x0 F. is suitable, with a temperature between about 120 and about 140xc2x0 F. being preferred.\nCalcium oxide is added to the fatty acid feedstock in the range from about 2 to about 3 equivalents relative to the fatty acid feedstock. A calcium oxide level between about 2.25 and 2.75 equivalents is preferred.\nWater is then added to hydrate the calcium oxide to its hydroxide form, creating a large amount of exothermic heat. The heat that is evolved usually is sufficient for the fatty acid neutralization reaction to proceed to completion, so that it is not necessary to supply heat to the reaction mixture. Between about 2 and about 5 equivalents of water relative to the calcium oxide is added to the reaction mixture, with between about 2.5 and about 3.5 equivalents being preferred.\nThe excess water is converted to steam by the exothermic heat generated and boils off rapidly. The reaction can be performed under atmospheric pressure, or under vacuum to draw off the steam.\nThe amount of time required for the reaction is between about 5 and about 60 minutes, and more typically between about 6 and about 15 minutes. The reaction is easily identified by the transformation of the admixture into a solid granular mass. Upon transfer from the reaction vessel, it can easily be processed into dry free-flowing particles.\nFish oil derived fatty acid calcium salt rumen bypass feed supplements of the present invention have a residual glyceride content below about 5 weight percent and contain from about 1 to about 10% by weight EPA calcium salt and from about 1 to about 10% by weight DHA calcium salt. Products containing from about 1.5 to about 9% by weight of EPA and DHA are preferred. The EPA- and DHA-containing feed supplements derived from fish oil/PFAD blends enhance the fertility of female ruminants. The present invention makes possible the commercially viable preparation of DHA- and EPA-rich calcium salt feed supplements for the enhancement of female ruminant fertility.\nThe present invention therefore includes fatty acid calcium salts having a fatty acid profile consistent with a profile resulting from blending from about 15 to about 50% by weight of fish oil with from about 85 to about 50% by weight of PFAD, wherein the fatty acid profile includes from about 1 to about 10% by weight of DHA and from about 1 to about 10% by weight of EPA.\nThe fatty acid calcium salt rumen inert feed supplements of the present invention may be conveniently fed to a ruminant admixed with a conventional ruminant feed. The feeds are typically vegetable materials edible by ruminants, such as legume hay, grass hay, corn silage, grass silage, legume silage, corn grain, oats, barley, distiller\"\"s grain, brewer\"\"s grain, soya bean meal and cottonseed meal. There is no particular lower limit of the calcium salt to be added to the ruminant feed, although in practice, amounts of the calcium salt below an amount that supplies 6 grams per day each of DHA and EPA are too small to provide significant fertility enhancement.\nThe fertility of female ruminants is enhanced when they are fed the EPA- and DHA-containing calcium salts starting as early as 21 days before parturition. While useful enhancement of fertility may be obtained by starting dietary supplementation at parturition or as late as 28 days following parturition, optimal results are obtained by earlier supplementation.\nThe dietary supplementation should continue daily at least until conception. However, because the DHA- and EPA-containing calcium salts of the present invention also enhance fertility by reducing embryonic death in the months following conception, the present invention also includes methods in which dietary supplementation is continued for at least 30 days and preferably for at least 60 days after conception. Beneficial results are obtained when dietary supplementation is continued up to 150 days after conception.\nThe DHA- and EPA-containing calcium salts of the present invention may be co-administered with additional quantities of other fatty acid calcium salts that are fed to ruminants for other purposes. The DHA- and EPA-containing calcium salts may be co-administered with a fatty acid calcium salt intended to supply energy to a high milk production ruminant, with a calcium salt having a fatty acid profile selected to modify the fatty acid profile of the milk fat or meat fat of the ruminant, or with both. Accordingly, the methods of the present invention for enhancing female ruminant fertility do not exclude the co-administration of other fatty acid calcium salt that do not contain DHA or EPA. One method in accordance with the present invention transitions the ruminant from the DHA- and EPA-containing calcium salts of the present invention to fatty acid calcium salts intended to supply energy to a high milk production ruminant once the fertility-enhancing benefits of the DHA- and EPA-containing calcium salts are no longer needed.\nThe following non-limiting examples set forth herein below illustrate certain aspects of the invention. All parts and percentages are by weight unless otherwise noted, and all temperatures are in degrees Fahrenheit."}
{"text": "1. Technical Field\nThe present invention relates generally to an improved data processing system and in particular to a method and apparatus for processing data. Still more particularly, the present invention provides a method and apparatus for transferring data on a bus between master and slave components.\n2. Description of Related Art\nA bus is a common pathway, or channel, between multiple devices. The computer's internal bus is known as the local bus, or processor bus. This type of bus provides a parallel data transfer path between the CPU and main memory and to the peripheral buses. A 16-bit bus transfers two bytes at a time over 16 wires; a 32-bit bus uses 32 wires, etc. The bus is comprised of two parts: the address bus and the data bus. Addresses are sent over the address bus to signal a memory location, and the data is transferred over the data bus to that location.\nVarious other types of buses are used in data processing systems. In particular, an inter-internal control (IIC) bus is an example of another type of bus used in a data processing system. In this bus, one wire carries a clock signal, while another wire carries the data signal. This type of bus is used to provide interconnection between various devices, such as a flexible service processor, a memory, and a control panel. A flexible service processor is a processing unit that is used to initialize a data processing system. A flexible service processor has its own boot code and operating system and may be connected to a number of I/O devices.\nDevices connected to an IIC bus may operate in a master or slave mode. In a master read/write mode, the device may initiate a data transfer. When a device is in a slave read/write mode, the device simply waits for data coming over the IIC bus. On a flexible service processor, a device driver is employed by the flexible service processor to handle data transfers with the bus. An application executing on a flexible service processor may make a call to the device driver indicating that the flexible service processor needs to respond as a slave to some event triggered by an external master. An external master is a master device located external to the flexible service processor. An example of an external master device is a control panel, which may send identifying button pushes on the panel, to the flexible service processor. Another example of an external master device is a rack power controller, which may send data, such as temperature data to the flexible service processor.\nMore than one master, external to the flexible service processor, may target a “virtual slave” in the flexible service processor. With this situation, more than one event may be issued by different external devices at the same time. Further, more than one process executing on the flexible service processor may be waiting for data on the same bus. A flexible service processor is the slave device and is only able to listen for data on a single slave address for a particular IIC bus. IIC buses have no intelligence regarding the identity of the master triggering the event. As a result, the bus is unable to associate data, put on data lines, with the process waiting for the data.\nWith this set-up, only one slave transfer request is currently possible for a device driver on a flexible service processor. If data for events are issued by different master devices, data can be lost if the flexible service processor is not continuously set up to receive any data that might be put onto the bus. For example, two applications may be executing on a flexible service processor in which one application monitors for button pushes on a control panel, while the other application monitors for temperature data from a power device, such as a rack power controller. If both the control panel and the rack power controller send data, as master devices, on the bus at the same time, the data is directed towards the flexible service processor as the slave device. This data can be received by only one of the applications because the current architecture only allows the device driver for the flexible service processor to handle only one slave request at a time. In this example, both applications cannot issue slave requests that can be handled by the device driver. As a result, a loss of information may occur.\nTherefore, it would be advantageous to have an improved method, apparatus, and computer instructions for handling data transfers between external master devices and internal requests for slave transfers."}
{"text": "1. Field of the Invention\nThe present invention concerns protective devices for pipes that have threading.\n2. Background of the Invention\nPipes, particularly pipes for oil fields, are often exposed to various climatic conditions when transported and stored. As a result pipes tend to corrode, primarily on their inside. Usually pipes with threading are provided with protective caps to protect the threads. Also, in order to prevent foreign bodies and rain water from reaching the interior of the pipes, closed or sealed protective caps and protective plugs are often used. But these caps and plugs result in the formation of condensed water due to variations in temperature. This water enhances the formation of rust.\nChemical agents that absorb moisture from their surroundings are well known in the prior art. Such agents are generally known as dessicants or drying agents. It is also known that in order to protect pipes from rust, dessicants may be packed in bags and inserted just inside the ends of pipes. These known protective devices do not suffice to protect pipes from the formation of rust, however. The moisture in pipes, particularly those of considerable length, is only absorbed by the dessicant that is in the immediate vicinity of the moisture. Thus, the inner pipe surfaces may rust despite the dessicant at the ends of the pipe. Thus with prior protective devices it is necessary to inspect the pipes from time to time by removing the protective caps and protective plugs. This visual inspection is quite cumbersome, time-consuming, and expensive."}
{"text": "The present invention relates generally to the quantification of microstructural fineness of metal castings, and more particularly to the automated quantification of dendrite arm spacing (DAS) in dendritic microstructures of metal castings as a way to avoid having to take such measurements manually.\nThe resulting microstructure of all cast aluminum-based components (such as engine blocks, cylinder heads, transmission parts or the like) is generally determined by the alloy composition and more particularly by the solidification conditions. In hypoeutectic alloys (i.e., those that contain less of the other alloying constituents than corresponds to the eutectic composition, examples of which include but are not limited to A356 and 319), the materials tend to solidify dendritically. Other such aluminum alloy examples that exemplify dendritic solidification include 354, 355, 360, 380, 383 and others. A typical microstructure of this family of alloys consists of a primary dendritic phase and a second phase of particles such as silicon particles and iron-rich intermetallics. The relative amounts, sizes and morphology of these phases in the as-cast structure are highly dependent on the solidification condition as well as on the alloy composition. The dendrite cell size (DCS) and DAS, sometimes referred to as the secondary dendrite arm spacing (SDAS), have long been used to quantify the fineness of the casting, which in turn can be used to gain a better understanding of the material and its related properties where—as a general rule—cast components with smaller DAS tend to have better ductility and related mechanical properties. Discussions pertaining to aluminum alloy casting in general, as well as to DAS properties in particular, may be found in numerous other applications for patent that are owned by the Assignee of the present invention, including U.S. patent application Ser. No. 12/356,226 filed Jan. 20, 2009, U.S. patent application Ser. No. 12/402,538 filed Mar. 12, 2009, U.S. patent application Ser. No. 12/454,087 filed May 12, 2009 and U.S. patent application Ser. No. 12/932,858 filed Mar. 8, 2011, all of which are hereby incorporated by reference.\nThere have been numerous efforts at describing dendrite refinement and its relationship to the solidification condition, starting in 1950 with Alexander and Rhines, who first established a quantitative basis for the influence of composition and solidification rate on certain dendrite features. Table 1 below summarizes known literature to describe in quantitative terms the fineness of dendritic structure.\nTABLE 1Microstructural parameters describing the dendritesParametersSymbolUnitDefinitionDendrite arm spacingDAS, λμmDistance between well(Levy et al., 1969;defined secondaryOswalt and Misra, 1980;dendrite arms (centerRadhakrishna et al., 1980;to center)Flemings et al., 1991)Dendrite cell sizeDCSliμmRandomly linear intercept(Spear and Gardner, 1963;among dendrite cellsJaquet and Hotz, 1992)Dendrite cell sizeDCSedμmArea equivalent circle(Cáceres et al. 1995)diameter of dendritecells including eutecticDendrite cell countCPUANumber of cells per field(McLellan, 1982)\nOf these, Spear and Gardner (1963) quantitatively described the scale of dendritic structure using DCS obtained by a random linear intercept and is referred to as DCSli in their FIG. 3(a). Following Spear and Gardner, Jaquet and Hotz (1992) in their study also used DCSli to quantify the dendrites. Levy et al. (1969), Oswalt and Misra (1980), Radhakrishna et al. (1980) and Flemings et al. (1991) all discussed DAS to quantify the dendritic structure. In these approaches, DAS is obtained by a linear intercept method where the line is chosen to intersect a series of well-defined secondary dendrite arms.\nMcLellan (1982) used dendrite cell count (CPUA) to quantify the microstructure and claimed that it describes the deformation process more accurately than DAS. However, Levy et al. (1969) had critically analyzed the measurements of both DAS and CPUA to characterize the cast structure, and pointed out that the standard deviation for DAS measurement was less than for CPUA measurement and also the mean cell size calculated from CPUA is greater than the mean DAS. Measurement of CPUA involves primary, secondary, and tertiary arms of the dendrites, whereas DAS measurements usually refer only to the secondary arm spacing.\nThe methods associated with manual measurement of DAS have been frequently used by the Assignee of the present invention as a way to make DAS measurement of aluminum castings. Such a procedure generally first includes preparation of metallographic samples that are prepared in accordance with known standards, such as the American Society of Testing and Materials Standard Guide for Preparation of Metallographic Specimens (also known as ASTM E3), a portion of which is reproduced below in Table 2.\nTABLE 2ASTM E3Abrasive Type/SizeForceAPlatenSurfaceLubricantANSI (FEPA)Time sec.N(lbf)RPMBRotationPlanar Grindingwater120-320 (P120-400)15-4520-30 (5-8) 200-300CCODpaper/stonegrit SiC/Al2O3Fine Grindingcompatible6-15 μm diamond180-30020-30 (5-8)100-150COheavy nylon clothlubricantRough Polishingcompatible3-6 μm diamond120-30020-30 (5-8)100-150COlow/no nap clothlubricantFinal Polishingcompatible1 μm diamond 80-12010-20 (3-5)100-150COmed/high nap clothlubricantsynthetic suedeEwater0.04 μm colloidal silica30-6010-20 (3-5)100-150CONTRAFor 0.05 μm aluminaAForce per 30 mm (1¼ in.) diameter mount.BPower heads generally rotate between 25 and 150 rpm.CHigh-speed stone grinders generally rotate at greater than 1000 rpm.DComplimentary rotation, surface and specimen rotate in same direction.EOptional step.FContra rotation, surface and specimen rotate in opposite directions.\nThe surface of the sample to be analyzed is expected to be of sufficient quality to reflect the truest possible size and shape of particles. In one form, the plane of the polish will include eutectic phases that will appear darker compared to the surrounding matrix. Thus, in one form, the metallographic samples are finally polished to obtain a flat, near mirror image surface finish. Chemical etching may be used to enhance the contrast of the dendrite structure, where in one form, the etching may be in accordance with ASTM E407. Preferably, the sample is clean and dry, while polishing artifacts (such as comet tailing, pitting, scratching, pullout and staining) should be kept to a minimum. Likewise, test conditions and deviations should be agreed upon beforehand. In a preferred form, each sample will be examined in numerous fields of view, each of which is subject to a high (for example, 100×) magnification, depending on the fineness of the material grain. After this, an image of the field of view to be measured should be captured. In one form, the linear intercept method may be used for measuring DAS, where three or more dendrites with visible dendrite trunks per field of view with at least three dendrite arms are selected. From this, a line is drawn from the outside edge of the first dendrite arm to the inside edge of the last dendrite arm; an example of this is depicted in FIG. 6B. The distance d for each dendrite may be recorded, while the number n1, n2, n3, etc. of dendrite arms counted for each measurement may also be recorded. These activities may be repeated for each field of view.\nAt present, both a volume percentage of eutectics and the DCS can be determined automatically via the use of an image analyzer. The local cooling rate affects not only the microstructural fineness but also porosity formation. Therefore, DAS tends to be used more frequently to quantify the microstructural fineness. The problem with the measurement of DAS is that it has to be performed manually by identifying well-defined dendrite arms in the image. Unfortunately, this is both very time-consuming and heavily dependent upon the skills of the user or individual doing the measuring."}
{"text": "This invention relates to space heating equipment generally, but more particularly to a flue gas heat recovery system for use in conjunction with a forced hot air home furnace. Most such furnaces are grossly inefficient, primarily because of heat lost to the atmosphere through the flue. This constitutes a continuous loss of heat, but it occurs at an accelerated rate when the burner is in operation, at which time flue gas temperatures are at a maximum, ranging from 350 to 400 degrees Fahrenheit.\nFully cognizant of this situation, the applicant has developed a heat recovery system which reduces the flue gas discharge temperature to approximately 125 degrees Fahrenheit. As a consequence of this substantial decrease in the discharge temperature, fuel consumption is significantly reduced, and the operating efficiency of the furnace greatly improved. The applicant's heat recovery system thus conserves energy and appreciably reduces home heating costs."}
{"text": "Flares are pyrotechnic devices designed to emit intense electromagnetic radiation at wavelengths in the visible region (i.e., light), the infrared region (i.e., heat), or both, of the electromagnetic radiation spectrum without exploding or producing an explosion. Conventionally, flares have been used for signaling, illumination, and defensive countermeasure decoys in both civilian and military applications.\nDecoy flares produce electromagnetic radiation through the combustion of a primary pyrotechnic material that is conventionally referred to as the “grain” of the flare. Flares, including decoy flares, configured to emit light in a visible spectrum of light may include a grain that includes magnesium and fluoropolymer-based materials. Including or excluding certain metals or other materials in the primary pyrotechnic material may alter the peak emission wavelength emitted by the decoy flare.\nDecoy flares are one particular type of flare used in military applications for defensive countermeasures. Decoy flares emit intense electromagnetic radiation at wavelengths in the infrared region of the electromagnetic radiation spectrum and are designed to mimic the emission spectrum of the exhaust of a jet engine on an aircraft.\nMany conventional anti-aircraft heat-seeking missiles are designed to track and follow an aircraft by detecting the infrared radiation emitted from the jet engine or engines of the aircraft. As a defensive countermeasure, decoy flares are launched from an aircraft being pursued by a heat-seeking missile. When the aircraft detects that a heat-seeking missile is in pursuit of the aircraft, one or more decoy flares may be launched from the aircraft. The heat seeking missile may, thus, be “decoyed” into tracking and following the decoy flare instead of the aircraft.\nFIGS. 1 and 2 illustrate a conventional flare 10, such as a decoy flare. Such conventional flares 10 include an elongated grain 22 that is inserted into a casing 12. The casing 12 may have a first end 14, i.e., the aft end, from which an aft end 23A of the decoy flare 10 is ignited, and a second end 16, i.e., the forward end opposite from the aft end, from which the grain 22 is ejected upon ignition. The flare 10 also includes a weighted end cap 40 that is attached to a forward end 23B of the grain 22. In some flares 10, the weighted end cap 40 may include an elongated rod 42 that is configured to be inserted into an internal bore 44 within the grain 22 to attach the weighted end cap 40 to the grain 22.\nThe weighted end cap 40 may be formed of a metal, e.g., brass, and may have a weight in the range of 30 grams to 50 grams. The weighted end cap 40 may be formed as a solid (i.e., monolithic) structure. The weighted end cap 40 is relatively small sized, compared to the grain 22, and relatively more dense than the grain 22. The majority of the weight of the grain 22 and weighted end cap 40, combined, is therefore distributed toward the forward end 23B of the grain 22 and, thus, of the flare 10. By locating the center of gravity toward the second end 16 of the flare 10, the weighted end cap 40 is configured to provide a stable trajectory to the flare 10 once ejected from the casing 12. In use, once the grain 22 of the flare 10 is combusted, the weighted end cap 40 remains essentially intact and falls to the ground below. The weighted end cap 40 falls below and behind the aircraft from which the flare 10 is ejected and, thus, presents danger to other airborne aircraft and ground-based building, vehicles, and personnel. Since multiple flares 10 may be fired at once during many evasive maneuvers, multiple weighted end caps 40 may form a so-called “cloud” of debris that is a danger to aircraft, building, vehicles, and personnel below."}
{"text": "The specificity of cytotoxic agents can be greatly improved by targeted delivery through linkage of the cytotoxic agents to cell-binding agents.\nMany reports have appeared on the attempted specific targeting of tumor cells with monoclonal antibody-drug conjugates (Sela et al, in Immunoconjugates 189-216 (C. Vogel, ed. 1987); Ghose et al, in Targeted Drugs 1-22 (E. Goldberg, ed. 1983); Diener et al, in Antibody mediated delivery systems 1-23 (J. Rodwell, ed. 1988); Pietersz et al, in Antibody mediated delivery systems 25-53 (J. Rodwell, ed. 1988); Bumol et al, in Antibody mediated delivery systems 55-79 (J. Rodwell, ed. 1988). All references and patents cited herein are incorporated by reference.\nCytotoxic drugs such as methotrexate, daunorubicin, doxorubicin, vincristine, vinblastine, melphalan, mitomycin C, and chlorambucil have been conjugated to a variety of murine monoclonal antibodies. In some cases, the drug molecules were linked to the antibody molecules through an intermediary carrier molecule such as serum albumin (Garnett et al, Cancer Res. 46: 2412 (1986); Ohkawa et al, Cancer Immunol. Immunother. 23: 86 (1986); Endo et al, Cancer Res. 47: 1076-1080 (1980)), dextran (Hurwitz et al, Appl. Biochem. 2: 25-35 (1980); Manabi et al, Biochem. Pharmacol. 34: 289-291 (1985); Dillman et al, Cancer Res. 46: 4886-4891 (1986); Shoval et al, Proc. Natl. Acad. Sci. 85: 8276-8280 (1988)), or polyglutamic acid (Tsukada et al, J. Natl. Canc. Inst. 73: 721-729 (1984); Kato et al, J. Med. Chem. 27: 1602-1607 (1984); Tsukada et al, Br. J. Cancer 52: 111-116 (1985)).\nA wide array of linker technologies has been employed for the preparation of such immunoconjugates and both cleavable and non-cleavable linkers have been investigated. In most cases, the full cytotoxic potential of the drugs could only be observed, however, if the drug molecules could be released from the conjugates in unmodified form at the target site.\nOne of the cleavable linkers that has been employed for the preparation of antibody-drug conjugates is an acid-labile linker based on cis-aconitic acid that takes advantage of the acidic environment of different intracellular compartments such as the endosomes encountered during receptor mediated endocytosis and the lysosomes. Shen and Ryser introduced this method for the preparation of conjugates of daunorubicin with macromolecular carriers (Biochem. Biophys. Res. Commun. 102: 1048-1054 (1981)). Yang and Reisfeld used the same technique to conjugate daunorubicin to an anti-melanoma antibody (J. Natl. Canc. Inst. 80: 1154-1159 (1988)). Dillman et al. also used an acid-labile linker in a similar fashion to prepare conjugates of daunorubicin with an anti-T cell antibody (Cancer Res. 48: 6097-6102 (1988)).\nAn alternative approach, explored by Trouet et al, involved linking daunorubicin to an antibody via a peptide spacer arm (Proc. Natl. Acad. Sci. 79: 626-629 (1982)). This was done under the premise that free drug could be released from such a conjugate by the action of lysosomal peptidases. In vitro cytotoxicity tests, however, have revealed that antibody-drug conjugates rarely achieved the same cytotoxic potency as the free unconjugated drugs. This suggested that mechanisms by which drug molecules are released from the antibodies are very inefficient.\nIn the area of immunotoxins, conjugates formed via disulfide bridges between monoclonal antibodies and catalytically active protein toxins were shown to be more cytotoxic than conjugates containing other linkers. See, Lambert et al, J. Biol. Chem. 260: 12035-12041 (1985); Lambert et al, in Immunotoxins 175-209 (A. Frankel, ed. 1988); Ghetie et al, Cancer Res. 48: 2610-2617 (1988). This was attributed to the high intracellular concentration of glutathione contributing to the efficient cleavage of the disulfide bond between an antibody molecule and a toxin. Despite this, there are only a few reported examples of the use of disulfide bridges for the preparation of conjugates between drugs and macromolecules. (Shen et al, J. Biol. Chem. 260: 10905-10908 (1985)) described the conversion of methotrexate into a mercaptoethylamide derivative followed by conjugation with poly-D-lysine via a disulfide bond. Another report described the preparation of a conjugate of the trisulfide containing toxic drug calicheamycin with an antibody (Hinman et al, Cancer Res. 53: 3336-3342 (1993)).\nOne reason for the lack of disulfide linked antibody-drug conjugates is the unavailability of cytotoxic drugs possessing a sulfur atom containing moiety that can be used readily to link the drug to an antibody via a disulfide bridge. Furthermore, chemical modification of existing drugs is difficult without diminishing their cytotoxic potential.\nAnother major drawback with existing antibody-drug conjugates is their inability to deliver a sufficient concentration of drug to the target site because of the limited number of targeted antigens and the relatively moderate cytotoxicity of cancerostatic drugs like methotrexate, daunorubicin, and vincristine. In order to achieve significant cytotoxicity, linkage of a large number of drug molecules, either directly to the antibody or through a polymeric carrier molecule, becomes necessary. However, such heavily modified antibodies often display impaired binding to the target antigen and fast in vivo clearance from the blood stream.\nIn spite of the above-described difficulties, useful cytotoxic agents comprising cell-binding moieties and the group of cytotoxic drugs known as maytansinoids have been reported (U.S. Pat. No. 5,208,020, U.S. Pat. No. 5,416,064, and R. V. J. Chari, Advanced Drug Delivery Reviews 31: 89-104 (1998)). Similarly, useful cytotoxic agents comprising cell-binding moieties and analogues and derivatives of the potent antitumor antibiotic CC-1065 have also been reported (U.S. Pat. No. 5,475,092 and U.S. Pat. No. 5,585,499).\nIt has also been shown that the linkage of highly cytotoxic drugs to antibodies using a cleavable link, such as a disulfide bond, ensures the release of fully active drug inside cells, and such conjugates are cytotoxic in an antigen specific manner (R. V. J. Chari et al, Cancer Res. 52: 127-131 (1992); U.S. Pat. No. 5,475,092; and U.S. Pat. No. 5,416,064).\nTaxanes are a family of compounds that includes paclitaxel (Taxol), a cytotoxic natural product, and docetaxel (Taxotere), a semi-synthetic derivative (see FIGS. 1 and 4), two compounds that are widely used in the treatment of cancer, E. Baloglu and D. G. I. Kingston, J. Nat. Prod. 62: 1448-1472 (1999). Taxanes are mitotic spindle poisons that inhibit the depolymerization of tubulin, resulting in cell death. While docetaxel and paclitaxel are useful agents in the treatment of cancer, their antitumor activity is limited because of their non-specific toxicity towards normal cells. Further, compounds like paclitaxel and docetaxel themselves are not sufficiently potent to be used in conjugates of cell-binding agents.\nRecently, a few new docetaxel analogs with greater potency than either docetaxel or paclitaxel have been described (I. Ojima et al, J. Med. Chem., 39: 3889-3896 (1996)). However, these compounds lack a suitable functionality that allows linkage via a cleavable bond to cell-binding agents (FIG. 1).\nThe synthesis of novel taxanes that retain high cytotoxicity and that can be effectively linked to cell-binding agents has been described recently (U.S. Pat. No. 6,340,701, U.S. Pat. No. 6,372,738 and U.S. Pat. No. 6,436,931, and FIGS. 2 and 4). In these disclosures, taxanes were modified with chemical moieties, ones containing thiol or disulfide groups in particular, to which appropriate cell-binding agents could be linked. As a result, these novel taxanes preserved, and in some cases even enhanced, the cytotoxic potency of known taxanes.\nIn the taxanes described in the aforementioned patents, the linking group was introduced at the C-10, C-7 or the C-2′ position of the taxane.\nIn the cases where the linking group was at C-7, the C-10 position did not have a free hydroxyl substituent but was rather an ester, ether or carbamate substituent. It has been previously shown (I. Ojima et al, J. Med. Chem., 39: 3889-3896 (1996)) that the presence of an ester or carbamate substituent at C-10 produced taxoids of high potency. However, there have been no studies on the potency of taxanes bearing a free hydroxyl group at C-10 and a linking group at C-7.\nAs described herein, the potency of taxanes bearing a free hydroxyl group at C-10 and a linking group at C-7 was found to meet or exceed the potency of taxanes bearing an ester, ether or carbamate substituent at C-10 and a linking group at C-7. Thus, in a first aspect, the present invention provides these novel taxanes bearing a free hydroxyl group at C-10 and a linking group at C-7 and having potent cytotoxic activity.\nFurther, in the taxanes described in the aforementioned patents, the linking group was introduced at the C-10, C-7 or the C-2′ position of the taxane. In all taxanes, the substituents at C-3′N and C-3′, were named —NHCOR4 and R3, respectively. Furthermore, the substituent at C-3′N, —NHCOR4, was either a benzamido group (R4=phenyl), like in paclitaxel, or a tert-butyloxycarbonylamino moiety (—NH-t-BOC, R4=t-butoxy), like in docetaxel. Based on published data, it was assumed that altering these substituents would cause a loss in potency. The substituent at C-3′ (R3) was either aryl or a linear, branched or cyclic alkyl group having 1 to 10 carbon atoms. Since, based on published data, it was thought that the substituents at C-3′N and C-3′ could not be altered without a loss in drug activity, the linking group was always introduced at a different position of the taxane, namely at C-7, C-10 or C-2′. Also, the inability to change the substituents at C-3′ or C-3′N greatly limited the variety of disulfide-containing taxanes that could be synthesized.\nIn a second aspect, the present invention is also based on the unexpected finding that the substituents at both C-3′ and C-3′N do not have to be limited to that present in the known taxanes. As described herein, the potency of taxanes bearing a variety of different substituents at C-3′N meets or exceeds the potency of taxanes bearing a benzamido or —NH-t-BOC substituent at this position. The new substituent at C-3′ or C-3′N can also contain a linking group that allows for linkage to cell-binding agents. The present invention discloses these novel highly potent taxanes bearing a variety of different substituents at C-3′ and C-3′N. The linking group can now be incorporated at any one of the five positions: C-3′, C-3′N, C-10, C-7 or C-2′."}
{"text": "Transparent, oriented polyester films have been known for a long time. Films produced from an antimony-free polyester are likewise known. Roasting bags and sleeves made of polyester film are likewise known.\nFilm bags for the baking of bread, and in particular film bags and tubes intended for cooking an entire dish in the oven, have to withstand temperatures above 200° C. for more than one hour without mechanical failure (e.g. due to bursting as a consequence of embrittlement). Industrial practice is to use almost exclusively film bags and tubes made of polyethylene terephthalate (PET), because this material has high heat resistance.\nThe traditional catalysts for producing polyesters continue to be antimony compounds. Some antimony compounds can be hazardous to health in particular at relatively high concentrations and if exposure is frequent. The EU therefore prescribes a maximum permissible migration of antimony from a film into a food. The high usage temperatures in oven applications increase antimony migration and for this reason alone it is desirable to reduce antimony content. Studies by the Swiss Department of Health show that, in actual foods, the migration limit for antimony is exceeded, sometimes considerably, at oven temperatures as low as 180° C. (see M. Haldimann, A. Blanc and V. Dudler, in Food Additives and Contaminants, 2007; 24(8): 860-868). Indeed, there are discussions in the technical literature where much smaller traces of antimony from PET are considered to be potentially hazardous (e.g. W. Shotyk and M. Krachler in Environ. Sci. Technol., 2007, (5), pp. 1560-1563). For an application in direct contact with food at greatly elevated temperatures it would therefore be important to use an antimony-free film.\nHowever, an analysis carried out for the purposes of this invention, taking more than 10 different roasting film bags and sleeves from various producers revealed the presence of antimony in all of the instances studied, albeit with values reduced (from 50 to 110 ppm) in comparison with standard antimony-containing films (standard film mostly >150 ppm). Antimony is easily detected at above 10 ppm by commonly used methods of digestion and analysis, and migration of antimony out of a film is similarly detectable by selecting suitable conditions (high temperature >100° C. in foods).\nThe reasons for the use of antimony compounds as polycondensation catalyst may well be found in the low heat resistance of antimony-free films. The oven test (see test methods) was not passed by any of the commercially available films studied that had been produced exclusively with a titanium catalyst. Germanium-catalyzed films were even poorer than the titanium-catalyzed films.\nAlthough experiments with use of phosphorus stabilizers, as described in JP 2007077220, to improve performance in the oven test by deactivating the catalyst led to reduced embrittlement, they did not lead to any reliable passing of the oven test. Nor was it possible to achieve reliable passing of the test by reducing, or entirely eliminating, use of regrind. Again, this achieved only some degree of improvement. This type of reduction moreover considerably impairs the cost-effectiveness of film production. Nor did any specific selection of the transesterification catalyst (elements studied being Zn, Mg, and Mn) lead to roasting film bags and sleeves which reliably pass the oven test.\nEP-A-2 164 079 (whose United States equivalent is United States Patent Application Publication No. 2010/209722) describes an electrical insulation film made of a polyester which was produced at least to some extent by using titanium catalyst, where the film comprises a mixture of free-radical scavengers, in order to achieve greater thermal stability. However, the high content of free stabilizer, and a selection of the stabilizers that is possible within the context of the disclosure can lead to films which are unsuitable for the direct food-contact application (e.g. by exceeding limiting migration values). Necessary preconditions for a roasting film bag or sleeve are not only thermal stability and the suitability of all of the starting materials for direct contact with foods but also suitable transparency and color giving good visibility of the roasting contents for the user, and also minimized clouding and yellowing during roasting, in order to provide good-quality appearance, and also compliance with certain mechanical properties which permit reliable handling of the bag or sleeve prior to and after roasting."}
{"text": "Handheld electronic devices can include mobile telephones, personal digital assistants (PDAs), handheld computers, laptop computers, or similar devices that might be able to communicate wirelessly with a telecommunications network and/or with the Internet. Other handheld electronic devices could include portable media players, such as music players and video players. Combinations of such devices, such as mobile telephone/PDA combinations and mobile telephone/music player combinations are also available.\nMany such devices, such as mobile telephones and portable music players, can produce an audio output and might consist of a handset portion that sends digital audio information wirelessly to a headset portion. For example, a mobile telephone might include a base handset that can communicate with a telecommunications network and an earpiece that can send and receive voice-related information wirelessly to and from the handset. The handset portion of such a handheld wireless device will be referred to herein as a base or a base unit. The earpiece portion will be referred to herein as a headset.\nTransmission of digital audio information between a base and a headset might occur via the Bluetooth wireless communication protocol or a similar technology. For example, a mobile telephone user with a Bluetooth-enabled base and a Bluetooth-enabled headset could carry on a hands-free, wireless telephone conversation. The headset converts Bluetooth signals received wirelessly from the base into speech audible to the user and converts the users speech into Bluetooth signals that are transmitted wirelessly to the base. The base may be worn on a belt, placed in a pocket, or otherwise be maintained some distance from the headset. A user might be able to answer a phone call by pushing a button on the headset and might not need to have direct access to the base to answer the call. The communication protocol between the base and the headset wilt be referred to herein as Bluetooth, but it should be understood that other technologies, whether currently in existence or developed in the future, could be used to achieve similar results and should be considered within the scope of the present disclosure."}
{"text": "Infectious disease affect multiple organs systems and are responsible for significant morbidity, mortality, and economic impact. Infectious agents most often present as a complex polymicrobial infections rather than as a single pathogens. Within the body, these polymicrobial infections cooperate with each other through mutualism changing both the type of antibiotics the organisms are susceptible to but also the level of antibiotics required to treat the infection as well as the virulence of the individual pathogens. The current gold standard is culture and sensitivity; a process in which a sample of is placed on a plate containing an agar medium for an overnight incubation allowing for individual organisms to be isolated and identified. An additional day is usually required to identify which antibiotic the individual organism is susceptible to through antibiotic susceptibility testing. The effectiveness of the this process is limited by the inherently poor sensitivity and specificity of the methodology. For example, recent studies have found that up to 25% of all urine culture results produce a false negative finding and that up to ⅔ of all uropathogens are missed by culture. This process which typically takes 48-72 hours has a significant impact on clinical practice in that clinicians are forced to treat the patient empirically without an accurate diagnosis and without proper guidance regarding the type of antibiotic that would be effective for treating the infection.\nAn inherent flaw in current antibiotic susceptibility testing is that it relies on testing only an individual organism rather than the entire pool of organisms simultaneously. By necessitating the isolation it adds a minimum of one additional day to the testing process. There is a need for a test that allows for the simultaneous assessment of antibiotic resistance without the need for isolation allowing for a more rapid determination of antibiotic resistance. The method described herein is a simple cost effective methods for assessing antibiotic resistance in all constituents simultaneously without the need for prior isolation."}
{"text": "Individuals and enterprises are continually exposed to risk because of future events beyond their control. The outcome of those events can either positively or negatively impact on their wellbeing.\nIndividuals and enterprises should generally prefer not to face exposure to the possibility of adverse consequences, regardless of their perception of the likelihood of such events occurring. It is in their interest to consider foregoing `resources` they currently possess if doing so would reduce the possibility of being so greatly exposed to future outcomes.\nRisk can take many forms in view of the large range and type of future events which might result in adverse consequences. Risk can be categorised, in one instance, as `economic` in nature. Phenomena that constitute economic risk include: commodity prices, currency exchange rates, interest rates, property prices, share prices, inflation rates, company performance, and market event based indices.\nAnother characterisation of risk concerns `technical` phenomena. This can include things like the breakdown of an electricity generation plant, aircraft engine failure, and the damage to, or failure of, orbiting telecommunications satellites. The outcomes for each of these phenomena will be adverse for the users and/or supplier.\nOther forms of risk defy ready characterisation, such as weather-based (viz., rain damage or lightning strike), or other natural occurrences (viz., earthquakes or iceberg collision with sea-going vessels).\nThere are also less tangible risks associated with, for example, the emission of atmospheric pollutants or the disposal of intractable toxic wastes, in the sense that the future consequences are unknown, save that there is a notion, based on current information, that they could be adverse.\nThe capability to manage risk is more important today than it was in the past, and is likely to become ever-more important into the future, because there is an ever increasing exposure to a wider generic range of future phenomena beyond the control of individuals or enterprises. There is also a wider feasible range of possible future events, and greater uncertainty about the likelihood of occurrence, associated with any single future phenomenon viz., an increasing volatility.\nIt is also thought that individuals are now more risk-averse in recessionary times, when there are fewer available discretionary resources to trade-off to protect themselves from such adverse future events.\nIn the prior art, individuals and enterprises faced with `technical` risk have hedged against future outcomes by mechanisms such as the adoption of quality assurance practices, warranties, increased research and development activity (and associated intellectual property rights such as patents, utility models and registered designs), the purchase of modernised plant and equipment, and improved inventory, occupational health and safety and employer/employee relations practices.\nConsider a manufacturer of, say, integrated circuits (ICs), which has many clients wishing to purchase its ICs. The demand may result in a delay in delivery due to limited manufacturing capacity, thereby requiring advance production scheduling for orders already in-hand. Typically, the manufacturer will give a warranty to a purchaser as to measurable performance criteria for its ICs; if a batch does not perform to the specified criteria, the manufacturer is required by contract to replace that batch. That is, a purchaser may have no interest in obtaining monetary compensation for the poor quality ICs, as the purchaser needs the components for their own products. In that case, the `consideration` the warranty makes is the priority scheduling of a substitute batch of that type of IC, possibly displacing other scheduled production runs, or deferring delivery to another purchaser.\nSuch contractual arrangements are piece-meal in nature, and can only be struck between the manufacturer and each individual purchaser. They also leave the manufacturer exposed to claims from other customers whose orders are delayed by the re-scheduling. The manufacturer has no convenient mechanism available to it to hedge against such claims, perhaps by way of reserving production rights with another manufacturer, in lieu of unavailability of their own manufacturing facility.\nIn the face of such `economic` risk, it is known for individuals and enterprises to hedge against adverse outcomes by indirect means such as self-insurance, and directly by means such as futures contracts, forward contracts, and swaps.\nThere are disadvantages or limitations associated with such available economic risk management mechanisms. Particularly, they provide, at best, only indirect approaches to dealing with the risk management needs. The available mechanisms are relatively expensive, and provide limited phenomenon coverage, and therefore cannot meet the requirements of the party seeking to hedge against such wide-ranging future risk. The infrastructure and pay-out costs associated with switching between, say, a commodities market and a stock market are often prohibitive for entities small and large alike. As a consequence, entities find themselves saddled with obligations they have little control over and cannot escape.\nIn respect of the \"less tangible\" forms of risk, an example in the prior art of a form of management of that risk is that of `pollution rights` sold by the U.S. Environmental Protection Agency (EPA) in March 1993 for the atmospheric emission of sulphur dioxide. This was done by an auction of \"allowances\" permitting the release into the atmosphere. By the year 1995, any company or organisation emitting sulphur dioxide in the U.S. without enough allowances to cover their total emissions will face prosecution. This means polluters must either buy further allowances, or else modify or replace their plant and equipment to reduce these emissions. The EPA will regulate the total number of allowances able to be obtained. The existing allowances have already become a valuable tradeable `property` as between sulphur dioxide emitters, that is, even before the time when no further allowances will be able to be purchased.\nManagement techniques for the `less tangible` forms of risk are in their infancy. The existing forms indicate an emerging demand for systems and methods to enable effective management.\nSpecific examples in the prior art of patents relating to methods and apparatus which deal with various forms of risk management include British Patent No. 2 180 380, in the name of Merrill Lynch Pierce Fenner and Smith Incorporated, directed to an Automated Securities Trading Apparatus (corresponding to U.S. Pat. No. 4,674,004, and further related to U.S. Pat. Nos. 4,346,442 and 4,376,978). Other examples include U.S. Pat. No. 4,739,478 assigned to Lazard Freres and Co., directed to Methods and Apparatus for Restructuring Debt Obligations, U.S. Pat. No. 4,751,640 assigned to Citibank, N. A., directed to An Automated Investment System, and U.S. Pat. Nos. 4,752,877, 4,722,055, and 4,839,804 assigned to College Savings Bank directed to Methods and Apparatus for Funding Future Liability of Uncertain Cost.\nThe present invention comes about in view of the shortcomings of existing risk management mechanisms, and the perceived increasing importance of the management of risk relating to specified, yet unknown, future events.\nIn this sense, the invention is directed to something having economic value to individuals, enterprises and societies as a whole. Methods and apparatus that provide for the management of risk offer material advantages by, for example, minimising adverse future outcomes, providing both a form of compensation in the event of adverse future outcomes, and forms of risk management not otherwise supported or available in the prior art, and thus have value in the field of economic endeavour."}
{"text": "1. Field of the Invention\nThe present invention relates to an improved integrated memory circuit, with the capability of reading or programming continuously, pages of data, without any gaps, and to methods of operating same.\n2. Description of the Prior Art\nSemiconductor integrated memory circuit devices for storing data typically have been categorized as either volatile, in which the data is lost once the power is turned off, or non-volatile, in which the data is retained even after the power is turned off.\nNon-volatile memories, comprising an array of non-volatile memory cells arranged in a plurality of rows and columns (or bit lines), can be categorized as either NAND or NOR type, referring to the manner in which the non-volatile memory cells are arranged in the array. Further, the non-volatile memory cells can be arranged to operate in a page mode manner, in which a page of data (typically 512 bytes) is stored in a plurality of latches (or plurality of page buffers) that are integrated with the memory circuit device. Reading of the integrated memory circuit device causes data from a page of the memory cells to be read and stored in the plurality of latches. Thereafter the contents of the plurality of latches are read, typically, in a serial manner, from the integrated memory circuit device. Programming of the integrated memory circuit device causes data from the external to be stored in the plurality of latches. Thereafter, the contents of the plurality of latches are stored in a page of non-volatile memory cells. Typically a page of non-volatile memory cells lie in the same row or word line.\nIn a conventional page-mode read operation, whenever a word line is addressed, a wait state is necessary for the on-chip control circuits to sense out data stored in the memory cells of that selected word line. After being sensed out, the data is latched into the plurality of page buffers before they are clocked out to the I/O pads. This wait-state, typically around several micro-seconds, accounts for a significant portion of the average page-mode read access time. Especially in applications of reading large volume data, several consecutive word lines are often addressed successively. With one wait-state for each addressed word line, the overall read performance is deteriorated. Therefore, a non-volatile memory with minimized number of wait-states is needed to provide high performance page-mode read operation.\nIn U.S. Pat. No. 5,768,215 a proposed solution to the aforementioned problem of the wait state is to provide two groups of page buffers, with each group of page buffers being one half the size of a page of memory cells. Initially, the data from a first page of memory cells is read into the two groups of page buffers. Thereafter, a first group of page buffers is read and the contents outputted to the external. However, as soon as the reading of the contents of the second group of page buffers commences, the reading of a second page of memory cells commences with the data read from one half of the second page of memory cells being stored in the first group of page buffers. After the contents of the second group of page buffers is outputted to the external, one half of the second page of memory cells will also have been read and stored in the first group of page buffers. As the reading of the first group of page buffers commences, the reading of the second half of the second page of memory cells commences and is stored in the second group of page buffers. This alternation of reading one half of a page of memory cells and storing the data into one of the groups of page buffers, while the contents of the other group of page buffers is read out continues.\nIn a conventional page-mode program operation, data is first loaded into the plurality of page buffers sequentially. Regardless of the number of bits (or collective bytes) to be programmed, all the data loaded into the plurality of page buffers will be programmed into the memory cells of a selected page simultaneously. Since the on-chip circuits, such as the charge pump (because typically programming requires a voltage source higher than the externally supplied voltage) can deliver only a limited amount of current, the efficiency of programming will deteriorate as the number of bits increases. Because of the limitation in the amount of current that can be provided by the on-board charge pump, one solution is to require a larger amount of time to program a page (or more) of data. Therefore a new technique is needed to provide high efficiency programming.\nFinally, non-volatile memory cells used in NAND architecture are typically of the stack gate type, such as that disclosed in U.S. Pat. No. 5,768,215. Further, the non-volatile memory cells used in NOR architecture can be both the stack gate type or the split gate type such as that disclosed in U.S. Pat. No. 5,668,757, whose disclosure is incorporated by reference in its entirety.\nIn the present invention, an integrated circuit memory device has a page of memory cells arranged in a plurality of sub-pages of memory cells. The memory cells are electrically coupled to a plurality of word lines, and a plurality of bit lines. A plurality of sub-page buffers are electrically coupled to the plurality of bit lines for storing data read from the memory cells coupled to the plurality of bit lines. The device further has an I/O data buffer. Each sub-page comprises a plurality of non-adjacent bit lines with memory cells coupled thereto, with the bit lines of each of the sub-page interleaving bit lines of another sub-page. A read controller circuit is coupled to the plurality of sub-page buffers for initiating a read operation to read data from a first sub-page of memory cells to an associated first sub-page buffer, while simultaneously reading data from a second sub-page buffer to said I/O data buffer.\nThe present invention also relates to an integrated circuit memory device having a NOR architecture emulating the read and programming operations of a NAND integrated circuit memory device. Finally, the present invention relates to a read-modify-write circuit."}
{"text": "In a general embodiment, the invention relates to an exercise and flexibility apparatus that may keep the body in good shape. In a preferred embodiment, the invention relates to a golf exercise and flexibility apparatus, and particularly to golf exercise apparatus which provides resistance to a golfer during a golf swing to strengthen and condition the muscles of the axial skeleton of the golfer.\nWhile this invention is described in terms of exercise and golf, the device may be used broadly for general conditioning, physical therapy and other sports such as swimming, tennis and the like where conditioning and flexibility are desirable. In one embodiment, this invention helps to prevent or minimize minor muscle aches and pains.\nCurrently, golf is an activity enjoyed by many people of all ages possessing varying degrees of athletic ability, musculoskeletal strength, flexibility and endurance. Although it is possible to perform a golf swing without having excessive musculoskeletal support, greater bodily strength, flexibility and endurance allows a golfer to hit a golf ball farther and with greater accuracy and consistency and to minimize minor muscle related aches and pains.\nExternal devices are currently being marketed to help train the muscles of the golfer to move along a predetermined path which is thought to be along an optimal golf swing path. These devices restrict the swing path of the golfer to a plane within which it is thought necessary to maintain the golf club throughout the golf swing. However, no resistance is supplied in the direction of rotation of the shoulders and upper torso, the hips, and upper legs of the golfer during performance. These devices are not designed to benefit muscular conditioning or flexibility.\nRegular exercise may keep the body in good shape, but not all exercise is equally effective. Many exercise devices on the market, particularly in health and athletic clubs, are less effective than patrons may assume, particularly as related to the rotational movements required in golf and similar activities. The problem is that most available equipment in health clubs train in predominantly linear, single plane movement and are limited to isolating one muscle group. The body rarely moves in just one plane and often requires multiple muscle groups to work together. Most body movement involves rotation and diagonal patterns of movement.\nThe need exists for an exercise, conditioning and rehabilitation device which permits activity consisting of components of motion in all three planes, and permits isolation of a specific area of the body, the motion of which is most desired. Such a device will permit a physical therapist, chiropractor or trainer or other instructional devices to tailor the activity of the user to the goals of the user."}
{"text": "The present invention relates to a reflective liquid crystal display using thin film transistors and a method of forming the same.\nA conventional reflective liquid crystal display using thin film transistors disclosed in the Japanese laid-open patent publication No. 5-224186 will be described with reference to FIG. 1. The reflective liquid crystal display is formed on an insulation substrate 1 such as a glass substrate. A gate electrode 2 is selectively formed on the insulation substrate 1. A gate insulation film is provided which extends over the gate electrode 2 and the top surface of the insulation substrate 1. A semiconductor layer 4 is selectively provided over a part of the gate insulation film 3 extending over the gate electrode 2 so that the semiconductor layer 4 covers over the gate electrode 2. A drain electrode 6 is selectively provided on a part of the semiconductor layer 4. A source electrode is provided on other part of the semiconductor layer 4 but separately from the drain electrode 6 so as to form an opening between the source and drain electrodes whereby a center part of the semiconductor layer 4 is shown by the opening for allowing a light to be incident through the opening into the semiconductor layer 4. The semiconductor layer 4 may be made of an amorphous silicon or a polysilicon so that when the semiconductor layer 4 absorbs a light having been incident, this generates carriers in the semiconductor layer 4. The semiconductor layer 4 may optionally be doped heavily with an n-type impurity to form an ohmic contact with the drain electrode 6. A reflective plate 5 serving as a pixel electrode is formed which extends over the gate insulation film 3. The reflective plate 5 is unitary formed with the source electrode which is positioned over the other part of the semiconductor layer 4. A color filter 10 is provided which extends over the reflective plate 5. An opposite electrode 8 is provided to be spaced apart from the reflective plate 5. The opposite electrode 8 is arranged which extends in parallel to the insulation substrate 1. The opposite electrode 8 is supported on a bottom surface of an opposite substrate 9. A liquid crystal is filled in the space defined between the opposite electrode and the reflective plate 5. The reflective plate 5 is made of a highly reflective material, for example, aluminum. The transistor formed over the substrate is an. inverted stagger transistor.\nOperations of the above transistor will subsequently be described. A gate voltage is applied to the gate electrode 2, then the drain electrode 6 is made conductive through the semiconductor layer 4 to the source electrode unitary formed with the pixel electrode in the form of the reflective plate 5. As a bias is applied across the drain electrode 6 and the pixel electrode, a current flows through the semiconductor layer 4 between the drain electrode 6 and the pixel electrode until the drain electrode 6 is the same in potential as the source electrode unitary formed with the pixel electrode in the form of the reflective plate 5. After the application of the gate voltage Vo the gate electrode 2 is discontinued, no current flows between the drain electrode 6 and the pixel electrode. Charges remain stored in the pixel electrode. As a result, a voltage between the opposite electrode 8 and the pixel electrode is controllable by the transistor. The liquid crystal is optically changed by the variation in voltage between the drain electrode 6 and the pixel electrode. The color filter 10 allows a color display. The opposite electrode 8 and the opposite substrate 9 are made of an optically transparent material.\nIncident light is transmitted through the opposite substrate 9 and the opposite electrode 8 into the liquid crystal 7. The incident light is then transmitted through the color filter 10 to reach the surface of the reflective plate 5 and reflected by the reflective plate 5. The reflected light is then transmitted through the color filter 10 and the liquid crystal 7 and further transmitted through the opposite electrode 8 and the opposite substrate 9. Other incident light is transmitted through the opposite substrate 9 and the opposite electrode 8 into the liquid crystal 7 and then transmitted through the opening between the drain electrode 6 and the source electrode to the semiconductor layer 4. The incident light is irradiated onto the semiconductor layer 4 and absorbed by the semiconductor layer 4. The absorption of the incident light causes generation of carriers, for example, electron-hole pairs in the semiconductor layer 4. Even if the pixel electrode was charged to apply an electric field to the liquid crystal 7, then the generation of carriers, for example, electron-hole pairs in the semiconductor layer 4 causes a current flowing between the pixel electrode and the drain electrode whereby the pixel electrode is discharged. As a result, the intensity of the electric field applied to the liquid crystal is dropped. This means all of the pixel electrodes are not charged and a reduced or almost no electric field is applied to the entire of the liquid crystal 7. This further means that contrast of the liquid crystal display is lowered.\nIt is, therefore, required to prevent the light from incidence into the semiconductor layer 4 or to shield the semiconductor layer 4 from the incidence of the light.\nAs described above, the reflective plate 5 is made of aluminum which has a high reflectivity but not chromium having a low reflectivity because if the reflective plate 5 is made of aluminum, then a sufficiently bright display can be obtained. If, however, the reflective plate 5 is made of chromium, then a dark display is obtained, which is unavailable.. The reflective plate 5 made of aluminum however has a disadvantage in diffusion of aluminum into the semiconductor layer 4 made of polysilicon or amorphous silicon whereby the characteristic of the transistor is deteriorated.\nFurther more, the semiconductor layer 4 is made contact through a thin orientation film such as a polyimide film with the liquid crystal 7. The polyimide film is incapable of blocking ions in the liquid crystal 7 for allowing the ions passing through the polyimide orientation film to the semiconductor layer 4, for which reason the semiconductor layer 4 is contaminated with the ions in the liquid crystal 7. As a result, the characteristic of the transistor is deteriorated.\nIf, in order to settle the above problem, it was proposed to provide an inorganic insulation film over the semiconductor layer 4 for having the inorganic insulation film cover the semiconductor layer 4. The inorganic insulation film may comprise a silicon nitride film having been formed by a plasma enhanced chemical vapor deposition method. The inorganic insulation film cover the semiconductor layer 4 is required to selectively be removed for allowing opposite side regions to be shown so that source and drain electrodes may be formed on the opposite side regions. The insulation film may selectively be removed by, for example, photolithography but this additional process causes increase in the manufacturing cost of the reflective liquid crystal display.\nIn the above circumstances, it had been required to develop a novel reflective liquid crystal display free from the above problems in deterioration in display caused by generation of carriers, for example, electron-hole pairs in a semiconductor layer connecting the source/drain electrode due to incidence of light into the semiconductor layer.\nAccordingly, it is an object of the present invention to provide a novel reflective liquid crystal display free from the above problems.\nIt is a further object of the present invention to provide a novel reflective liquid crystal display free from deterioration in display caused by generation of carriers, for example, electron-hole pairs in a semiconductor layer connecting the source/drain electrode due to incidence of light into the semiconductor layer.\nIt is a still further object of the present invention to provide a novel reflective liquid crystal display which may be formed at a relatively low manufacturing cost.\nIt is yet a further object of the present invention to provide a novel reflective liquid crystal display capable of a bright display with high contrast.\nIt is moreover object of the present invention to provide a novel reflective liquid crystal display including thin film transistors showing highly reliable operations.\nIt is further more object of the present invention to provide a novel reflective liquid crystal display having thin film transistors showing highly reliable operations.\nIt is further more object of the present invention to provide a novel reflective liquid crystal display having a reflective plate with a sufficiently large thickness and a rough surface but free from excess reflection.\nIt is still more object of the present invention to provide a novel reflective liquid crystal display with a wide visible angle.\nIt is yet more object of the present invention to provide a novel reflective liquid crystal display having gate electrodes with a reduced resistance but free of any substantial delay in transmission of signals.\nIt is a still further more object of the present invention to provide a novel reflective liquid crystal display free from deterioration in characteristics of a thin film transistor in the reflective liquid crystal display.\nIt is yet further more object of the present invention to provide a novel reflective liquid crystal display having thin film transistors having semiconductor layers which are free from exposure to liquid crystal in the reflective liquid crystal display.\nIt is an additional object of the present invention to provide a novel reflective liquid crystal display having a semiconductor layer free from deterioration due to contamination to the semiconductor layer with aluminum caused by contacting an aluminum reflective plate to the semiconductor layer.\nIt is a further additional object of the present invention to provide a novel reflective liquid crystal display having an aluminum reflective plate free from deformations such as hillock.\nIt is another object of the present invention to provide a novel reflective liquid crystal display which may be formed without carrying out excess or additional fabrication processes at a relatively low manufacturing cost.\nThe above and other objects, features and advantages of the present invention will be apparent from the following descriptions.\nThe present invention provides a reflective liquid crystal display comprising the following elements. Source/drain electrodes are formed on an insulation substrate. The source/drain electrodes are separated from each other. A pixel electrode is provided which has an end portion unitary formed with the source electrode and the pixel electrode extending over the insulation substrate. A semiconductor layer extends over the insulation substrate in contact directly with a gap between the source/drain electrodes and also extends on at least respective parts of the source/drain electrodes for connecting the source/drain electrodes. A gate insulation film extends over entire parts of the semiconductor layer. A gate electrode extends on at least a part of the gate insulation film so that the gate electrode is positioned to cover the semiconductor layer and to be separated by the gate insulation film from the semiconductor layer for shielding the semiconductor layer from incident light. A reflective plate is formed over the pixel electrode so as to be separated from the gate electrode and from the semiconductor layer. An opposite electrode is supported on a surface of an opposite substrate so that the opposite electrode is separated by a space filled with a liquid crystal from the gate electrode and from the reflective plate and also that the opposite electrode faces to the gate electrode and the reflective plate."}
{"text": "In the attempts to understand human disease and develop therapies, the medical and associated fields have frequently relied upon analogies to observations in animals. While there are a number of similarities between non-primates and primates, there are also a large number of differences. Many of the naturally-produced agents, such as cytokines, hormones, surface membrane proteins, blood proteins, and the like, differ significantly in their composition. Most importantly, for the transduction of signals, cellular interactions, regeneration and differentiation, the molecules involved frequently have a narrow range of host specificity. Where autocrine or paracrine mechanisms are not involved, xenogeneic tissue introduced into a non-primate host must totally rely upon the host for the source of ligands not produced by the xenogeneic tissue or other xenogeneic tissue must be supplied which makes such ligands available. Furthermore, xenogeneic tissue is subject to rejection by an immunocompetent host, so that means must be provided to protect the xenogeneic tissue from degeneration as a result of immune rejection. There are no viable hosts which have no protective mechanisms against foreign entities.\nThere has been substantial interest in using immunosuppressive agents or immunocompromised hosts to extend the period in which xenogeneic explants may be maintained as viable functional tissue. Frequently, in order to be useful to study the effect of the agents on the xenogeneic tissue, the tissue must retain most if not substantially all of its natural characteristics. However, since the xenogeneic tissue, at best, is only a very small amount of tissue compared to the entire host, it can be subject to many modifications which may change its characteristics, so as to substantially reduce its predictive capability as to the response of the xenogeneic host.\nWhere the xenogeneic tissue has mobile cells, such as host organs of peripheral blood lymphocytes, such as lymph nodes, thymus, spleen, pancreas, tonsils and the like, these cells may be rapidly displaced by the xenogeneic host cells. Furthermore, vascularization and lymphatic connection can result in the xenogeneic tissue being invaded by host endothelial cells which provide the vessels. Other host cells may also invade the tissue, so as to change the characteristics of the tissue. It is therefore not sufficient that the tissue remain viable or even grow, it is necessary that the tissue retain a substantial degree of similarity to the native tissue, so as to provide responses which are predictive of the responses of the native tissue in the source host for the xenogeneic tissue.\nHumans as a host are particularly unique, in the many restrictions which apply to humans and the manner in which they may be studied. Normally, one cannot induce a disease in a human to study the etiology of the disease. Nor can one treat a human with an agent, without first having gone through extensive tests with other mammals, to determine efficacy and pharmacological properties. With many diseases, there are many infectious agents which are tropic for human tissue. In other pathologic conditions, human tissue may respond quite differently from other hosts, due to the nature of its surface membrane proteins, particularly as to metabolic pathways, pharmacology, and the like.\nIt is therefore of substantial interest that mammalian hosts, particularly small mammalian hosts, be developed, which allow for the study of the response of human tissue to a variety of agents, particularly pathogenic agents and therapeutic agents used to treat the infection or diseased state."}
{"text": "For many years, film cameras were the only option for capturing cinema quality motion pictures. The time requirements and costs related to shooting and processing motion picture images on film stock and then transferring those images into a digital form have created a need for motion picture cameras that capture high definition or cinema resolution imagery directly in a digital form. The advent of Digital Cinemas, cost effective Stereo 3D Digital cinema projection systems and establishment of Digital Cinema Initiative SMPTE Standards has fueled the need for more content creation for delivery at 2K, 4K and Stereo formats.\nAccordingly, there is a need for a digital camera system that meets the needs described above that reduces costs. There is also a need for a digital camera system that leverages digital processing and visualization tools. There is a further need for a digital camera system that provides user feedback and metadata collection when shooting special effects, compositions and stereo or multi-camera content. There is an additional need for a digital camera system that improves flexibility in networked collaboration, enables separated imaging block and recording, has a simple workflow with metadata management and, importantly, maintains film-like reproduction qualities and camera operation. There is also a need for a digital camera system that not only provides the capabilities described above but can also utilize existing cabling infrastructure, broadcast signal monitoring and transmission systems. There is a need for a digital camera system that mixes broadcast standard digital sources into a recording and visualization system, as well as generate broadcast standard and network streaming outputs for 2D and 3D content.\nIn the past few years, while several digital cinema cameras have emerged on the market these digital cinema cameras are complex designs with limited connectivity that are only able to address a limited set of the needs described above. For example, these digital cinema cameras are incompatible with existing cable infrastructure. Also, these digital cinema cameras either completely lack network management or are capable of only minimal network management (i.e., only simple controls that lack full image streaming or metadata management). Further, these digital cinema cameras lack the ability to capture or record multi-sensor 2K or 4K image data using a single control application. Additionally, these digital cinema cameras lack visualization tools or metadata integration. These digital cinema cameras do not utilize existing broadcast infrastructure to transmit multi-resolution data and have complex workflows with respect to stereo 3D and multi-camera content acquisition, broadcast and network transmission either live or in a post production process. These digital cinema cameras are limited to 1920×1080 image sensor pixel arrays that require the use of a multiple sensor prism block which, in turn, requires use of complex and expensive optics. These digital cinema cameras utilize dedicated hardware functions with no or limited image processing flexibility or upgrade capability. Dedicated hardware functions utilized by these digital cinema cameras include video processing to perform non-reversible color space transformations or sub-sampling to formats, such as YUV 4:2:2 and 4:4:4, as standard broadcast signals. These digital cinema cameras implement a variety of proprietary compression and coding schemes that introduce visible image artifacts, especially when projected on large screens. While a number of these digital cinema cameras can generate preview imagery for display on an electronic viewfinder, these digital cinema cameras can only do so with limited resolution or visualization tools. High-resolution outputs from these digital cinema cameras are restricted to transmission in SMPTE standard resolution and formats. These digital cinema cameras often output imagery to be record on restrictive, proprietary or large external storage devices. These storage devices include a tape storage system having only linear data access, Non-Volatile Flash or RAM drives with limited storage, and multiple hard disk drive RAID storage systems which are often non-portable and whose media cannot be easily removed or transported for direct use in other systems. Also, the files stored on these storage devices have limited color correction, image processing or post production metadata integration.\nIn recent years, many digital still cameras or dual-mode video and still camcorders have also been developed which use single image sensors with color filter arrays. These digital still cameras and camcorder devices do use higher resolution sensors (e.g., HD (1920×1080) camcorders, digital single-lens reflex camera (DSLR) are now 10 MP and higher). However, these digital still cameras and camcorders have slow readout architectures (e.g., a DSLR may only shoot four (4) frames per second) and can only achieve video rate preview at low resolution (e.g., 640×480) or standard definition (e.g., VGA 640×480 at thirty (30) frames per second) using sub-sampling or windowing techniques. These digital still cameras and camcorders use dedicated hardware functions or targeted function digital signal processors (DSP) to perform image processing to interpolate and colorize the raw image data from the image sensor. These digital still cameras and camcorders compress the colorized images for storage; but the compressing process performed by these devices prevents access to the original full raw image pixel data for later processing, analysis or colorization. In addition, the interpolation and color processing applied to the source raw data in those devices initially generates data sets that are larger than the source raw data which, in turn, requires the application of higher compression to fit the data sets into a target storage capacity. This typically results in a reduction in image quality compared to the original image or a coded version of the raw data.\nA few single sensor cameras have been developed for use in 2K and 4K acquisitions in raw format. However, these cameras use dedicated hardware or targeted function DSPs to perform image processing to interpolate, colorize and display preview quality output and simultaneously compress the raw sensor image data for later digital editing and grading. Also, the compression method and metadata employed by these cameras foreclose the dynamic retrieval of alternative resolution or representations at different bit rates during recording for network streaming, remote grading or adaptive editing. Due to their architectures, these single sensor cameras must apply high compression to fit data into target internal storage capacity devices. Also, due to their architectures, these single sensor cameras lack the ability to transmit the imager raw content over existing broadcast or network infrastructure cabling for remote monitoring, networking, recording or visualization. These single sensor cameras cannot process captured signals with prerecorded content or broadcast format signals for live composition, switching, grading, mixing into virtual sets or adding graphic overlays based on extracted metadata or analytics. These single sensor cameras also lack the ability to manage, control or record multi-sensor imagers, which may be remotely connected to a recording system.\nIn recent years, there has been an interest in producing digital cinema quality 3D stereographic, wide-dynamic and immersive content using multiple imagers. This has created a need for more efficient modular and scalable cameras and workflow solutions. There is a further need for a digital camera system having a precise synchronization mechanism to enable images to be mixed or stitched without motion artifacts. While digital camera systems have been used to produce this type of content, these camera systems suffer from the same system limitations as the cameras described above. These camera systems are mostly comprised of stand-alone cameras, each with individual controls, viewing and recording systems, with no integration mechanism other than a common sync signal (i.e., there is no communication between camera controls or viewing and recording settings). These camera systems are large and bulky such that the camera systems cannot be placed very close together physically, as is required for short inter-ocular distances in 3D stereographic or for creating hemispherical views where cameras need to be placed as close together as possible from a common center point. When shooting thru mirrors and beam splitters, rigs (i.e., a combination of digital cameras, optics and mounting platform) become more cumbersome and difficult to use in handheld shooting environments. Finally, these camera systems lack a comprehensive set of image processing, visualization, positioning control, recording, playback, communications and display tools for use in such high-definition multi-camera systems."}
{"text": "1. Field of the Invention\nThe invention generally relates to MICR documents and, in particular, to a method for tracking the source of the MICR encoded documents.\n2. Description of Related Art\nMagnetic ink character recognition (MICR) codes are commonly used on financial documents such as checks, loan coupons, travelers checks, bonds and the like. The MICR codes are printed on the checks by various methods. Currently, the majority of the checks printed are done on litho or offset presses with liquid ink which has iron oxide mixed into it. The usage for these checks is predominantly checkbooks, preauthorized payment checks, and travelers checks, along with other applications like payroll checks. Within the last few years, however, there has been a growing trend in the use of computer-driven impact and nonimpact continuous form and page printers to print the MICR codes, as well as other data on checks and financial documents. Although not limited to these computer-driven printers, it is mainly in this area that this invention will be used. Impact printers mostly use a carbon film ribbon with iron oxide mixed in it to print with. Nonimpact page and continuous form printers use toner or thermal ribbon with iron oxide mixed in it to print MICR codes. The iron oxide (ferrite) is magnetized. When the MICR code is read utilizing a magnetic reader, a magnetic wave form or signature is produced which identifies the particular code being read. The MICR codes facilitate automatic reading and processing of the documents by automatic readers. Such readers can read up to 100,000 documents per hour. The font of the MICR codes is controlled by standards which cannot be deviated from.\nOccasionally, automatic MICR readers have difficulty reading the MICR code of a particular document. Most often, the rejected document has been produced on a non-impact printer which is either out of adjustment or is using a poorly designed or programmed MICR font. Unfortunanely, conventional MICR documents include little or no information allowing one to determine the source of the MICR code placed on the document to thereby allow the defective MICR printer to be identified and corrected. The conventional MICR encoded documents may provide some information regarding the bank or other business organization which initially authorized the document. However, such information is frequently insufficient to quickly and accurately determine the actual machine that printed the MICR code on the document. This problem becomes particularly intractable where the issuing bank or other business organization may have contracted any of a number of document printing companies to print the documents. Furthermore, it has become feasible to print MICR documents using a laser printer supplied with magnetic ink toner, rather than using a conventional dedicated MICR printer. MICR documents printed with such laser printers have a fairly high incidence of subsequent processing errors.\nOccasionally, fraudulent financial documents are created by persons using NIP (nonimpact) printers without proper authorization, or using stolen font cards, soft fonts, or MICR printers. Currently, no convenient method is available for determining whether any particular MICR document was printed by a stolen MICR printer or printed using a stolen MICR font card or MICR font. The ability to make such a determination would greatly facilitate the identification of counterfeit or unauthorized MICR documents. Because of the worldwide standards governing the MICR codes, no other indicia or marks of any kind have been imbedded in a MICR cell, as the MICR reader would attempt to read it if it were in the clear band and it would result in a reject. However, the subject invention assures that any marks, codes, etc., can be printed along with the MICR code in the cell readability."}
{"text": "Field of the Invention\nThis invention relates to systems and methods for selectively blocking access to particular Internet websites and pages. More specifically, embodiments of this invention relate to a flexible filtering system and method that provides administrators with several options for controlling end-user access to those sites.\nDescription of the Related Art\nThe Internet is a global system of computers that are linked together so that the various computers can communicate seamlessly with one another. Internet users access server computers in order to download and display informational pages. Once a server has been connected to the Internet, its informational pages can be displayed by virtually anyone having access to the Internet.\nThe easy access and inexpensive cost of retrieving Internet pages has led to several problems for controlling access to inappropriate information, such as pornography. Several solutions to this problem have been proposed, including rating systems similar to that used for rating movies so that a parent or employer could control access to Internet servers, or pages, that have a particular rating. Unfortunately, this mechanism requires each person running an Internet server to voluntarily rate their site. Because of the free-wheeling nature of the Internet, this type of voluntary rating scheme is unlikely to be very efficient for preventing access to sites, such as those containing pornography, that most parents or businesses desire to block.\nIn addition to a rating scheme, others have developed databases that contain the uniform resource locater (URL) address of sites to be blocked. These databases are integrated into network computer systems and Internet firewalls so that a person wishing access to the Internet first has their URL request matched against the database of blocked sites. Any URL found in the database cannot be accessed by the user. One such system is described in U.S. Pat. No. 5,678,041 to Baker et al. Unfortunately, such systems rely on the database of accessed sites to be complete. Because new servers are being added to the Internet on a daily basis, as well as current servers being updated with new information, these databases do not provide a complete list of sites that should be blocked.\nIn addition, current systems do not provide a user with any flexibility once they have requested a site that is within the blocked database. Thus, users that have a legitimate reason for reviewing such sites are still denied access.\nThus, what is needed in the art is a flexible system that provides control over which sites can be retrieved, but also has the flexibility to provide access to blocked sites under specific conditions. The present invention provides such a system."}
{"text": "One of the aims of future cellular wireless communication systems is to enhance the achievable data throughput to mobile terminals (MT) situated close to a cell edge. This is important as, assuming a reasonably uniform distribution of MTs over the cell area, then a significant fraction of the MTs in a cell is close to its periphery. When the same carrier frequency is re-used in neighboring cells, the signal from the ‘wanted’ base station (BS) with which the MT is communicating is received at the cell edge at power levels similar to signals originating from BSs in neighboring cells. Cell edge MTs therefore experience strong interference in addition to low signal to noise ratio (SNR), which makes it difficult to achieve high data rates to these MTs.\nIn current third generation (3G) systems like the Universal Mobile Telecommunications System (UMTS), macro-diversity and soft-handover techniques are known. These allow simultaneous communication between more than one BS and a MT in order to improve the link quality to MTs at the boundaries between cells. In macro-diversity and soft-handover the same data is transmitted to a MT from multiple BSs. This is achieved by having all transmissions effected on the same carrier frequency, and discriminating transmissions from different BSs based on their different scrambling codes. The MT comprises a receiver arrangement for receiving multiple (CDMA) signals simultaneously, i.e. it has multiple receive signal paths (descrambling and decorrelation) and a combiner to combine the despread symbol streams.\nMacro-diversity techniques in 3G systems rely upon code division multiple access (CDMA) techniques in order for a MT to receive a given data stream from more than one BS simultaneously. However, each wanted data stream causes interference to the reception of the other stream.\nFor future cellular systems, including 3rd Generation Partnership Project (3GPP) Long Term Evolution (LTE) and Wireless World Initiative New Radio (WINNER), multi-carrier (OFDM-based) transmission schemes are being proposed (at least for the downlink). Further, multiple access (sharing the time-frequency resources between MTs) is typically envisaged to be based on allocating different groups of subcarriers to different MTs (orthogonal frequency division multiple access (OFDMA)) rather than on CDMA. The macro-diversity techniques from 3G are therefore not directly applicable to these future systems.\nOne favored approach for improving cell-edge performance in these future OFDM-based cellular systems is to partition subcarriers between cells, wherein all cells may use all subcarriers at lower transmit powers. This gives coverage to the inner parts of the cell but does not reach the cell edges (and therefore does not cause interference to neighboring cells). For communication to MTs at the cell edges a BS then uses a subset of the total number of available subcarriers where the subsets are defined to be mutually exclusive with the subsets employed in neighboring cells. This prevents higher power transmissions to cell-edge MTs from causing high interference to cell-edge MTs in the neighboring cells.\nThis approach improves the inter-cell interference situation for cell-edge MTs at the expense of increasing the frequency re-use factor, which results in lower spectral efficiency than re-using all subcarriers in all parts of every cell (i.e., frequency re-use factor of one). It may also reduce the peak throughput to cell-edge users since only a subset of the total number of subcarriers are available for use.\nA straight forward extension of the 3G macro-diversity ideas to these OFDM systems would be to use the cell-edge subcarrier subsets of two or more neighboring cells to transmit to a MT. The MT would then receive the same data from multiple BSs (via different subsets of subcarriers) and can combine these to enhance the data reception quality. The downside of this approach is of course that this consumes resources (subcarriers) in two or more cells for the benefit of one MT. This is analogous to 3G macro-diversity, which requires resources (spreading codes) to be allocated in two or more cells for the benefit of one MT, and requires the MT to receive and combine two or more signals.\nA related piece of prior art is the operation of Single Frequency Networks (SFN), which are known in broadcast systems such as Digital Audio Broadcasting (DAB) and Digital Video Broadcasting (DVB). In these OFDM systems the same data signal is broadcast from all transmitters. In the regions approximately mid-way between two transmitters, the receiving terminal receives a super-position of the signals from both transmitters. This is equivalent to receiving the signal from a single source via the composite channel given by the summation of the two channels from each transmitter. With a suitably long guard interval the receiver in these OFDM systems can successfully receive the combined signal from the two sources with enhanced signal strength over reception from a single transmitter, without Inter-Symbol Interference (ISI), and without needing to be ‘aware’ that the signal originated from two separate sources.\nFIG. 1 shows a schematic block diagram of a transmitter with a space-time coder 20 adapted to receive an input signal 10 and to generate two transmission signals 30 which can be jointly received at a receiver.\nHowever, although the SFN concept means that a simple receiver can be used, the combined signal can still undergo fading.\nFor the case of two transmission paths with respective transfer functions h1 and h2 to the receiver antenna, the combined transfer function becomes (h1+h2), so that the SNR of the received signal is (h1+h2)2/n2 where n is the amplitude of noise and interference. But sometimes the particular values of h1 and h2 will cancel, significantly reducing the received power. Therefore the received signal quality could be highly variable."}
{"text": "Various orthopaedic surgical procedures involve the use of a number of orthopaedic instruments, including orthopaedic trials, for example. For example, during a typical total knee arthroplasty (TKA) tibial trials, such as tibial trial trays and tibial trial inserts, are used to assist a surgeon in preparing the tibial surface for implantation of the tibial portion of the artificial knee. A surgeon often uses a tibial insert trial to determine the tibial implant size and to ensure a proper alignment and tibial component thickness prior to implanting the tibial components themselves, for example.\nCertain types of orthopaedic instruments such as saw blades, for example, are made of metal, but others, such as trials for example are oftentimes formed from molded polymers. Many such orthopaedic instruments are cleaned and sterilized after use in a particular orthopaedic surgical procedure and then reused in subsequent surgical procedures. In a typical cleaning or sterilization process, the tools may be autoclaved. In such a process, the reusable orthopaedic tools are placed in sterilization trays for cleaning and sterilization."}
{"text": "To make retractable a very large roof so as to move it completely off the space it covers presents new challenges. Such a roof may be a barrel vaulted roof hundreds of feet long, or a pitched roof hundreds of feet long, or even a flat roof hundreds of feet long. The object is to move such a roof as one unit. The object also is to move the flat roof in sections.\nThe state of the art today shows no quick answer to these requirements. There are smaller retractable roofs often made of plastic or glass used over swimming pools and garden courts. There are unit skylights of various shapes either in glass or plastic that are made to slide. There are custom roofs in either glass, plastic, or other material. None of these roofs show how to move a large barrel vaulted roof, a pitched roof, or a flat roof in the large size required, completely on and off the opening they cover.\nWhen the roofs become very large, one of the new issues is the temperature expansion of the retractable roof versus the supporting structure. In new stadiums this is sometimes accounted for by partially articulated wheel trucks. On smaller roofs this is sometimes accounted for by expansion and contraction of the retractable roof frame taken by sideway movement of the wheels on the supporting rails of the roof. In this invention the effort is not to mitigate the expansion and contraction movement of the roof on the wheels, but to create a system where the rail for the roof supporting wheels moves closely the same as the roof from temperature expansion and contraction, so no thermal expansion and retraction accommodating means are needed at all between the roof and the wheels. This is done by supporting the rail on an intermediate frame which spans the area being covered so that the rail and the support therefore will expand and contract as one. To manage this, the intermediate frame is supported on slide bearings resting on columns, or possibly walls or another beam structure, so that it is free to move in unison with the roof, both exposed to similar temperatures. In this manner the roof and wheels see limited to small differential movement so no measures such as articulated frames or rollers or even slide bearings are needed between the roof and the wheels.\nThis alone does not assure that the roof can be moved evenly. The retractable frame that is moved must be very rigid. Rigidity in panelized roofs on stadiums extends only to the individual panels. This invention moves the entire roof as a monolithic piece and therefore requires bracing of the entire roof. This is done very simply with X bracing or other forms of bracing in the retractable roof.\nThirdly to move a large monolithic roof evenly requires that the drive means, the wheels with attached motors or a cable and winch system, the two most common, must work evenly together. Other stadium roofs use various means of electro mechanical controls to assure one side moves the same as the opposite side so the roof moves evenly. Some stadiums use controls that measure the exact location of the sliding wheels at all times and correct motor speeds to adjust continuously so that the roof runs evenly. In other words so that one side reaches the end point at the same time as the other. Although this could be done with the roof in this invention, this invention uses X bracing in a horizontal frame supporting the roof to transfer differential traction forces directly between the wheels and motors which drive the roof. Direct transfer of traction forces evens out the movement of all the motors and wheels contributing to an even parallel movement of the roof. This invention also uses motors on all wheels. At the same time redundancy is achieved. Should one motor fail, the loss will be taken up by the others.\nIt is these ideas brought together in an artful and engineered manner that result in a smooth operating very large retractable roof to be used on a mall or other large area, which has never been done in this or a similar manner.\nThe following drawings together with the detailed description will describe this further."}
{"text": "Rate matching is a technique used in mobile communication systems to match a code rate of an encoder such as a convolutional encoder or turbo encoder to the data transmission rate of a communication channel. Rate matching typically involves puncturing or repeating coded bits output by the encoder to match the data transmission rate of the communication channel. Rate matching allows a single encoder to be used for a plurality of data channels with different data transmission rates.\nIn a conventional rate-matching circuit, an encoder receives an input bit stream and generates two or more coded bit streams. An interleaver interleaves each coded bit stream. A rate matching circuit bit-multiplexes the coded bits in each interleaved bit stream and outputs a single bit stream to the transmitter having a desired number of bits to match the data transmission rate of the communication channel. The bit-multiplexing performed by the rate matching circuit intermixes the interleaved bits from all interleaved bit streams. If the number of bits output by the encoder is greater than the number required, some of the interleaved bits are punctured. Conversely, if the number of bits output by the encoder is less than required, some of the bits may be repeated. The rate matching circuit may be implemented using a circular buffer, or a real-time multiplexing circuit.\nWhile rate matching circuits used in the past provide good performance, there remains a need for new rate matching circuits for convolutional codes with lower complexity that provide good performance."}
{"text": "Electroplating is a well-known process used, for example, to nickel plate workpieces such as automotive parts that are required to have a \"bright\" finish. In a typical example of this particular application of the process, the workpiece passes successively through one or more nickel plating baths and associated rinse baths. In each case, the bath contains a nickel salt solution and a nickel anode. The workpiece acts as the cathode and is electroplated by connecting a source of direct electric current between the cathode and anode. In one particular plating system, the workpiece receives a thin nickel coating in a first bath (often called the \"nickel strike\"). Additional nickel is deposited on the workpiece in the second and third baths which are referred to respectively as the \"semi-bright\" and \"bright\" baths. The nickel salt solutions in these baths contain organic \"brightners\".\nIn a commercial nickel plating operation, plating is carried out on a \"batch\" basis. For example, a batch of automotive parts such as bumpers may be carried by a single rack by which they are transported from bath to bath. The bumpers are electrically connected together to effectively form a single cathode and are all electroplated at the same time. As the workpieces leave the plating baths a thin film of plating solution remains on the surface. This is referred to as \"dragout\". Some dragout occurs no matter what the shape of the workpiece, but the amount can be substantial with workpieces of unusual shapes. Liquid transfer into a bath in this way is referred to as \"dragin\".\nDragout from metal plating baths represents a significant cost in terms of the value of the lost metal as well as the cost of treatment of the water used to rinse the workpieces after they have been plated. Probably even more significant is the cost and difficulty of disposal of hazardous metal hydroxide waste sludge that conventional waste treatment systems produce.\nDespite the depletion effect of dragout from the plating baths, in some cases, it is unnecessary to add metal salt to the plating bath because the metal concentration in the bath tends to increase naturally due to poor cathode efficiency. In extreme cases, it may even be necessary to periodically decant some of the plating solution and replace it with water to prevent a build-up in the metal concentration in the plating bath."}
{"text": "Field of the Invention\nThe present disclosure relates to a display device. More particularly, the present disclosure relates to a display device having improved light-use efficiency.\nDescription of the Related Art\nA non-self-emissive display device, such as a liquid crystal display device, an electrophoretic display device, an electrowetting display device, etc., requires a separate backlight unit to generate light. The backlight unit is classified into an edge-illumination type backlight unit and a direct-illumination type backlight unit according to a position of a light source with respect to a display screen on which an image is displayed. The edge-illumination type backlight unit has a relatively thinner thickness than that of the direct-illumination type backlight unit. Accordingly, a mobile display device mainly employs the edge-illumination type backlight unit."}
{"text": "The present invention relates to an industrial mat for use in construction sites that require heavy lifting, crane support stabilization, temporary areas for installation of pipelines or bridges, and other support structures in locations where the ground requires grading and stabilization to provide a flat support surface for such equipment or activities.\nConventional industrial mats are typically 4 feet wide and utilize beams having square or rectangular cross sections of sizes between 8×8 inch to 16×16 inch that are up to 40 feet in length. Typically, the beams are made of oak and preferably white oak as that material provides acceptable performance of the mats for a significant service life at a reasonable cost. Such mats are available from Quality Mat Company, Beaumont, Tex. These mats, which are often called timber mats or crane mats, typically utilize virgin wood that is shaped and cut to length to meet design demands. Due to weather conditions and other environmental factors, however, the availability of trees that can be harvested to make such large size and length beams is reduced, thus making it difficult to obtain suitable quantities to make large numbers of mats.\nAccordingly, alternatives are needed for industrial mat constructions to conserve the amount of wood beams that need to be included. Also, the materials that may be considered as alternatives need to possess the necessary physical properties to be able to withstand harsh outdoor conditions as well as to support heavy equipment. And of course cost is a factor in determining the selection of alternate materials, as it is not cost-effective to provide a mat that is multiple times more expensive than one that can be made of wood.\nAdditionally, there are concerns regarding the lifting and movement of such mats. In the past, lifting, manipulation and placement of the mats were facilitated by exposing part of the joining member for grasping by a lifting hook of a crane or other heavy piece of equipment. While the exposure of the joining or connecting rod facilitated lifting or moving operations, it presented an issue in that workers could inadvertently step into over the open hole that was provided around the connecting rod. Thus, other designs are needed to provide more safely to personnel working on these mats.\nThe present invention now satisfies this need in the art for such alternative constructions by providing newly designed industrial mats that exhibit improved performance over total wood mats at a competitive cost. These new mats also include improved structures and arrangements for lifting and manipulation of the mats. The new mats also avoid the reliance upon long wood beams for the entire mat construction as in the past."}
{"text": "Technical Field\nThe present invention relates to a magnetically coupled reactor and a power converter including the magnetically coupled reactor. The magnetically coupled reactor includes a coupled core member having a first core and a second core. The magnetically coupled reactor also includes a first coil and a second coil that are wound around legs of the coupled core member.\nRelated Art\nThere are cases where it is desired to decrease leakage magnetic flux in a magnetically coupled reactor by increasing the degree of magnetic coupling. To this end, a technique such as disclosed in JP 2009-284647 A has been conceived.\nJP 2009-284647 A has disclosed a magnetically coupled reactor formed by two outer cores and an inner core disposed between the two outer cores. Two coils are wound around the outer cores and the inner core. One of the two coils is wound around one of the outer cores and the inner core, while the other coil is wound around the other of the two outer cores and the inner core. The coils are partly overlapped alternately to form a lap-winding portion, and the lap-winding portion is then disposed around the inner core.\nDuring the use of the magnetically coupled reactor, a negative-phase-sequence current that flows in a reverse direction may flow through the two coils. In this case, magnetic flux is generated in both coils and becomes intensified at portions where the two coils face each other, which therefore may cause an increase of Joule loss. If the lap-winding portion is formed as in the structure of JP 2009-284647 A, the Joule loss could be decreased. However, such a structure may require new cores at both ends in a direction orthogonal to the axial direction of the coils, and extended windings to be wound around the two cores also become necessary. As a result of this, the size of the structure is increased. Therefore, in a power converter, improvement is desired from the view point of miniaturizing the structure in which a voltage converting operation can be performed using the leakage inductance without excessively suppressing the leakage inductance."}
{"text": "This invention is in the field of semiconductor integrated circuits. Embodiments of this invention are more specifically directed to reliability screening of non-volatile storage elements in such integrated circuits.\nConventional metal-oxide-semiconductor (MOS) and complementary MOS (CMOS) logic and memory devices are prevalent in modern electronic devices and systems, as they provide an excellent combination of fast switching times and low power dissipation, along with their high density and suitability for large-scale integration. As is fundamental in the art, however, those devices are essentially volatile, in that logic and memory circuits constructed according to these technologies do not retain their data states upon removal of bias power. Especially in mobile and miniature systems, the ability to store memory and logic states in a non-volatile fashion is very desirable. As a result, various technologies for constructing non-volatile devices have been developed in recent years.\nA recently developed technology for realizing non-volatile solid-state memory devices involves the construction of capacitors in which the dielectric material is a polarizable ferroelectric material, such as lead zirconate titanate (PZT) or strontium-bismuth-tantalate (SBT). Hysteresis in the charge-vs.-voltage (Q-V) characteristic, based on the polarization state of the ferroelectric material, enables the non-volatile storage of binary states in those capacitors. In contrast, conventional MOS capacitors lose their stored charge on power-down of the device. It has been observed that ferroelectric capacitors can be constructed by processes that are largely compatible with modern CMOS integrated circuits, for example placing capacitors above the transistor level, between overlying levels of metal conductors.\nFIG. 1 illustrates an example of a Q-V characteristic of a conventional ferroelectric capacitor. As shown, the charge (Q) stored across the conductive plates depends on the voltage applied to the plates (V), and also on the recent history of that voltage. If the voltage V applied across the capacitor plates exceeds a “coercive” voltage +Vα, the capacitor polarizes into the “+1” state. According to this characteristic, once polarized to the “+1” state, so long as voltage V remains above coercive voltage −Vβ, the capacitor exhibits a stored charge of +Q1. Conversely, if the voltage V applied across the capacitor plates is more negative than coercive voltage −Vβ, the capacitor is polarized into the “−1” state, and will exhibit a stored charge of −Q2 for applied voltage V below +Vα.\nAn important characteristic of ferroelectric capacitors, for purposes of non-volatile storage in integrated circuits, is the different capacitance exhibited by a ferroelectric capacitor in its polarized states. As fundamental in the art, the capacitance of an element refers to the ratio of stored charge to applied voltage. In the context of a ferroelectric capacitor, the change in polarization state that occurs upon application of a polarizing voltage is reflected in charge storage. For example, referring to FIG. 1, the polarization of a ferroelectric capacitor from its “−1” state to its “+1” state is reflected in a relatively high capacitance C(−1), by way of which polarization charge involved in the change of polarization state is retained within the capacitor as the voltage exceeds its coercive voltage Vα; on the other hand, a capacitor already in its “+1” state exhibits little capacitance C(+1) due to polarization, since its ferroelectric domains are already aligned prior to the application of the voltage. In each case, the ferroelectric capacitor also has a linear capacitance, by virtue of its construction as parallel plates separated by a dielectric film (i.e., the ferroelectric material). As will be evident from the following description, a stored logic state is read by interrogating the capacitance of ferroelectric capacitors to discern its polarized state.\nFerroelectric technology is now utilized in on-volatile solid-state read/write memory devices. These memory devices, commonly referred to as “ferroelectric RAM”, or “FeRAM”, or “FRAM” devices, are now commonplace in many electronic systems, particularly portable electronic devices and systems. FeRAM memories are especially attractive in implantable medical devices, such as pacemakers and defibrillators.\nFIG. 2a illustrates one example of a conventional FeRAM memory cell, as described in Miwa et al., “A 512 Kbit low-voltage NV-SRAM with the size of a conventional SRAM”, Digest of Technical Papers: Symposium on VLSI Circuits, Paper 12-4 (IEEE, 2001), pp. 129-132, incorporated herein by reference. Cell 10 is constructed essentially in the form of a typical CMOS static RAM (SRAM) cell, with cross-coupled inverters forming a latch. In this example, one inverter is formed of p-channel transistor 2P and n-channel transistor 2N with their source-drain paths connected in series between power supply node PWR (which, in normal operation, provides the power supply voltage Vdd) and ground, and the other inverter is formed of p-channel transistor 3P and n-channel 3N, with their source-drain paths also connected in series between power supply node PWR and ground. These inverters are cross-coupled by storage node S1 at the drains of transistors 2P, 2N being connected to the gates of transistors 3P, 3N, and storage node S2 at the drains of transistors 3P, 3N being connected to the gates of transistors 2P, 2N. N-channel MOS pass transistor 4A has its source-drain path connected between storage node S1 and bit line BL, and its gate connected to word line WL; similarly, n-channel MOS pass transistor 4B has its source-drain path connected between storage node S2 and bit line BLB, and its gate also connected to word line WL. In a write operation, the complementary logic levels at storage nodes S1, L2 are written from a differential signal at bit lines BL, BLB upon word line WL turning on pass gates 4A, 4B; conversely, the state of storage nodes S1, L2 can be read by energizing word line WL and sensing the differential signal on bit line pair BL, BLB. The states at storage nodes S1, L2 remain latched while cell 10 is powered-up.\nFerroelectric capacitors 5A, 5B in cell 10 provide its non-volatile capability. In this example, capacitor 5A is connected between storage node S1 and plate line PL, and capacitor 5B is connected between storage node S2 and plate line PL. Once the state of cell 10 is present at storage nodes S1, L2, plate line PL remains at an intermediate level (e.g., Vdd/2) for both read and write operations, whether cell 10 is selected or not, so long as cell 10 is powered up. In this conventional case, capacitors 5A, 5B can be programmed to store a current state at storage nodes S1, L2 by driving plate line PL to a high voltage (e.g., power supply voltage Vdd) and then to a low voltage (e.g., ground). For example, if storage node S1 is at a high level and storage node S2 is at a low level, the high voltage level on plate line PL (greater than coercive voltage +Vα, considering plate line PL relative to storage node S2) will polarize capacitor 5B to a “+1” state (as shown in FIG. 1); capacitor 5A does not polarize at this point. Conversely, the low voltage level on plate line PL (more negative than coercive voltage −Vβ, considering plate line PL relative to storage node S1) will polarize capacitor 5A to a “−1” state; the polarization of capacitor 5B is not affected. These polarization states will remain on capacitors 5A, 5B, even after cell 10 is powered-down.\nTo recall the state of cell 10 as retained at capacitors 5A, 5B, plate line PL is held at a low voltage (e.g., ground) while cell 10 is powered up. Referring to FIG. 1, the powering-up of cell 10 effectively moves the applied voltage in the negative direction from the “0” voltage axis. As the voltage at power supply node PWR increases toward Vdd, the voltage across capacitors 5A, 5B crosses coercive voltage −Vβ, and capacitor 5B begins polarizing. The change in polarization of capacitor 5B causes its effective capacitance to be greater than that of capacitor 5A. This causes storage node S1 to rise toward power supply voltage Vdd faster than does storage node S2; the latch action of cell 10 will then drive storage nodes S1, S2 to their previous state, setting the state of the latch. The polarization state of capacitor 5B is destroyed in this recall (or “restore”) operation; this recall is thus a generally destructive read.\nFIG. 2b illustrates another example of a conventional FeRAM cell 10′, as described in Masui et al., “Design and Applications of Ferroelectric Nonvolatile SRAM and Flip-Flop with Unlimited Read/Program Cycles and Stable Recall”, Proceedings of the Custom Integrated Circuits Conference, Paper 16-6-1 (IEEE, 2003), pp. 403-406; and in Masui et al., “A Ferroelectric Memory-Based Secure Dynamically Programmable Gate Array”, J. Solid State Circuits, Vol. 38, No. 5 (IEEE, May 2003), pp. 715-725; both incorporated herein by reference. Cell 10′ of FIG. 2b is similarly constructed as cell 10 of FIG. 2a, with the same reference numerals used for similar elements, but includes two additional ferroelectric capacitors 7A, 7B. Capacitor 7A is connected between storage node S1 and plate line PL2, and capacitor 7B is connected between storage node S2 and plate line PL2; plate line PL1 refers to the plate line connected to capacitors 5A, 5B in this FIG. 2b. Capacitors 7A, 7B are much larger capacitors than their capacitors 5A, 5B, for example having three times the dielectric area and thus three times the linear capacitance. The operation of cell 10′ in its read and write operations is essentially the same as that of cell 10 of FIG. 2a, however both plate lines PL1, PL2 may be maintained at ground potential in normal operation for read and write cycles of cell 10′ as an SRAM cell. Programming of ferroelectric capacitors 5A, 5B, 7A, 7B in cell 10′ is accomplished in similar fashion as described above. In this case, both plate lines PL1, PL2 are driven to a high voltage (e.g., Vdd) and then to a low voltage (e.g., ground) while cell 10′ remains powered up. The logic level at storage node S1 determines the polarization of capacitors 5A, 7A, while the logic level at storage node S2 determines the polarization of capacitors 5B, 7B.\nPlate lines PL1, PL2 receive separate voltages during the recall operation in this cell 10′. Recall begins with plate lines PL1, PL2 at a low voltage (e.g., ground), while pass transistors 4A, 4B are held off by word line WL at an inactive low level. Plate line PL1 is then driven to a high voltage (e.g., Vdd) while plate line PL2 remains at a low voltage, at which point cell 10′ is powered up. With plate line PL1 at Vdd and plate line PL2 at ground, the voltages at storage nodes S1, S2 are based on a voltage divider of the capacitances of ferroelectric capacitors 5A, 5B, 7A, 7B, which depends on their differing polarization states.\nFor example, if storage node S1 is at a high logic level and storage node S2 at a low logic level, programming of ferroelectric capacitors 5A, 5B, 7A, 7B will result in ferroelectric capacitors 5A, 7A polarized to the “−1” state, and ferroelectric capacitors 5B, 7B polarized to the “+1” state. In the recall operation, with plate line PL1 at Vdd and plate line PL2 at ground, the voltage divider of ferroelectric capacitors 5A, 5B, 7A, 7B at each of storage nodes S1, S2 operates as:\n            V      s        =                  V        dd            ⁡              (                              C            5                                              C              5                        +                          C              7                                      )              ⁢        where VS is the voltage at one of storage nodes S1, S2, C5 is the capacitance of its associated ferroelectric capacitor 5, and C7 is the capacitance of its associated ferroelectric capacitor 7 (C7>C5 in this example, as described above). In this case, ferroelectric capacitor 5A at storage node S1 has a higher capacitance because of its “−1” polarization state than will ferroelectric capacitor 5B in its “+1” state. Accordingly, the voltage VS1 at storage node S1 will be higher than the voltage VS2 at storage node S2, by virtue of the higher effective capacitance of capacitor 5A than that of capacitor 5B. One can also consider this effect as storage node S1 being charged by the polarization charge from capacitor 5A as changes its polarization state; storage node S2 does not receive this polarization charge, because capacitor 5B is not changing state. As a result, as cell 10′ is powered-up, the differential voltage at storage nodes S1, S2 established by ferroelectric capacitors 5A, 5B, 7A, 7B is latched into the cross-coupled inverters of cell 10′, with storage node S1 latched high and storage node S2 latched low in this example. As before, this restore operation is destructive of the programmed state, requiring reprogramming of capacitors 5A, 5B, 7A, 7B when appropriate.\nFerroelectric capacitors are now also implemented in logic circuits, for example to store the state of latch nodes upon power-down, so that this previous state can be restored on power-up. Another use of ferroelectric capacitors in logic circuits is to provide non-volatile storage of configuration information, so that power-up of the logic circuit will immediately load that stored configuration information into various latch nodes within the logic circuit. These and other uses of non-volatile elements such as ferroelectric capacitors in logic circuits are becoming popular.\nFIG. 2c illustrates an example of conventional ferroelectric flip-flop 16, as used in the data or control flow of a large-scale logic circuit, and in which ferroelectric capacitors 15A, 15B, 17A, 17B operate to store the state of flip-flop 16 on power-down, and to restore that logic state as it is again powered-up. The construction and operation of an example of a flip-flop such as that shown in FIG. 2c is described in Masui et al., “Design and Applications of Ferroelectric Nonvolatile SRAM and Flip-Flop with Unlimited Read/Program Cycles and Stable Recall”, Proceedings of the Custom Integrated Circuits Conference, Paper 16-6-1 (IEEE, 2003), pp. 403-406, incorporated herein by reference. In the simplified example of FIG. 2c, input data line D is received at pass gate 9A, which connects input data line D to storage node S1 according to the level of clock signal CLK. Pass gate 9B selectively connects the output of inverter 8B to storage node S1. Pass gates 9A, 9B are each CMOS pass gates, constructed as a p-channel MOS transistor in parallel with an n-channel MOS transistor, with the gates of the two transistors receiving complementary versions of clock signal CLK (i.e., clock signal CLK itself and via inverter 11). In this example, pass gates 9A, 9B are complementary with one another, such that one is turned on when the other is turned off, and vice versa.\nInverter 8B is one of a pair of cross-coupled inverters providing the latch function in flip-flop 16, and receives storage node S2 at its input. The other inverter 8A in the cross-coupled pair has its input coupled to storage node 51, and drives storage node S2, completing the cross-coupling. In this case, the output of flip-flop 16 is taken on output data line Q from storage node S2, via inverter 14 (so that the output data level corresponds to that of input data line D). Bias to inverters 8A, 8B is provided from power supply node PWR by p-channel MOS transistor 13A and from ground by n-channel MOS transistor 13B. The gates of transistors 13A, 13B are driven by enable signal EN (inverted by inverter 19 in the case of transistor 13A).\nFerroelectric capacitors 15A, 15B, 17A, 17B in conventional flip-flop 16 provide the non-volatile retention of its state upon its power-down (e.g., enable signal EN driven inactive low). Similarly as in memory cell 10′ of FIG. 2b, ferroelectric capacitors 15A, 17A are connected between storage node S1 and plate lines PL1, PL2, respectively; similarly, ferroelectric capacitors 15B, 17BA are connected between storage node S2 and plate lines PL1, PL2, respectively. Ferroelectric capacitors 17A, 17B are typically larger than ferroelectric capacitors 15A, 15B, and provide the voltage divider load for the recall or restore operation, as described above in connection with FIG. 2b. The operation of flip-flop 16 in programming its ferroelectric capacitors 15A, 15B, 17A, 17B with the current state of flip-flop 16 (i.e., storage nodes S1, S2), and in restoring that state on power-up, follows that described above relative to FIG. 2b. \nCircuits such as memory cells 10, 10′, and flip-flop 16, provide important non-volatile functionality in many modern integrated circuits. Non-volatile retention of memory contents, and logic states, as provided by these circuits enable the use of high-performance integrated circuit functionality in many applications previously thought to be impractical in the absence of non-volatile storage capability. However, ferroelectric elements such as capacitors are known to be vulnerable to later-life degradation due to various effects. Repeated read/write cycles (i.e., polarization and destructive read operations) can degrade the polarizability of ferroelectric capacitors, as can the effects of time and temperature.\nFor example, the mechanism of “imprint” is a well-known “end-of-life” mechanism that has been observed in ferroelectric devices. Imprint is exhibited by permanent shifts in the shape of the Q-V characteristic loop, generally appearing as flattened slopes in the transition curves and shifting in the polarization charge levels that result in an overall loss of area within the hysteresis loop. One type of imprint, referred to as “dynamic” imprint, results from repetitive or extended application of voltage exceeding one of the coercive voltages Vα or Vβ, as the case may be, which effectively rewrites the same data in the cell. Imprint of the static type refers to permanent shifts in the polarization characteristic due to the capacitor remaining in one polarization state over time without external bias. FIG. 3 illustrates the effect of imprint on the Q-V characteristic of an example of a ferroelectric capacitor, by way of shifted characteristic curve 20 relative to its initial characteristic curve 22. In this somewhat extreme case, the effects of imprint include inadequate polarization of the capacitor in either state, as well as poorly defined coercive voltages.\nAs evident from FIG. 3, severe imprint greatly reduces the ability to distinguish the two polarization states from one another. In the context of memory cells 10, 10′ and flip-flop 16 described above, in which the polarization states restore the previous circuit state according to capacitance differences between the polarization states, one can see that the capacitances C(−1), C(+1) indicated by the slopes of the Q/V transitions for the two polarization states from zero applied volts to a positive voltage are much closer to one another after the effects of imprint have set in (curve 20), than at time-zero (curve 22). Because the restore or recall operation relies on differences in capacitance between the two states, this weakness in polarization increases the likelihood that device imbalances, or noise events, can cause the memory cell or logic element incorporating ferroelectric capacitors to fall victim to a data error on restoring a programmed state. Certainly the likelihood of such errors increases as degradation of the ferroelectric material over time, temperature, and usage occurs. Such data errors have been observed in modern high-performance integrated circuits incorporating ferroelectric capacitors for non-volatile data and logic state retention.\nIt has been observed that any sizable population of ferroelectric capacitors will exhibit a distribution of polarization characteristics, and thus will include some ferroelectric capacitors that are “weaker” than others upon manufacture, in regard to the ability to electrically distinguish between their polarization states. Under the reasonable assumption that degradation of the polarization capability over time, temperature, and usage is uniform across the population of elements, those capacitors exhibiting the weaker polarization characteristic at time-zero will be the first to cause circuit failure in use. It would therefore be useful to perform a time-zero (i.e., at wafer level, or after packaging, or both) electrical test to screen out, from a population of manufactured integrated circuits, those devices that will be the earliest to fail within a particular threshold time.\nOf course, in circuit implementations such as memory cells and logic circuits, it is difficult if not impossible to obtain a direct measurement of the polarization characteristics of each embedded ferroelectric capacitor. Rather, the polarization behavior of these elements can only be interrogated in an indirect manner. In the case of memory cells such as those described above in connection with FIGS. 2a and 2b, one conventional screening approach is to reduce the power supply (Vdd) bias applied to the memory cells 10, 10′ when performing a restore operation. This reduced power supply bias will be reflected in weaker distinction in the voltages presented by the ferroelectric capacitors to the storage nodes during the restore operation. However, it has been observed that there is a poor correlation between those devices failing such a reduced Vdd storage test, and those devices that exhibit the weakest reliability due to time-dependent polarization degradation, especially as reliability requirements continue to become more stringent.\nAnother conventional approach to time-zero screening of non-volatile latches in memory cells is to apply a small disturb differential voltage from the bit lines to the storage nodes, prior to performing a restore or recall operation. That approach is somewhat cumbersome insofar as construction of write circuits in the memories is concerned. In addition, this bit line disturb voltage approach is not applicable to logic circuits such as flip-flop 16 that do not have accessible “bit lines” at all, especially for such circuits that are embedded in complex logic chains and data paths.\nBy way of further background, copending and commonly assigned U.S. patent application Ser. No. 12/142,568, published as U.S. Patent Application Publication No. US 2009/0316469 A1, incorporated herein by reference, describes a test procedure for screening weak ferroelectric capacitors in a non-volatile memory, using a short high temperature bake."}
{"text": "In the sense of the method of the invention, the term \"fiber strands\" is used to include both ordered and irregular strands of textile single fibers (slivers) and furthermore also webs.\nIt is known in the textile industry to use mechanical or capacitive sensors to monitor automatically mass fluctuations of continuously moving fiber strands or slivers for purposes of quality control. On the other hand, foreign particles and foreign fibers are measured manually by a lab technician who randomly pulls apart slivers on a luminous table and visually searches for foreign particles and counts them. The properties of the individual fibers cannot be measured at all, and therefore it is impossible to ascertain changes in parameters (for instance length, crimping) caused by processing conditions."}
{"text": "Fluid transport systems are well known and used in a number of applications. One specific application of transporting a fluid in a machine is the transportation of ink in a printer. Common examples of inks include aqueous inks and phase change or solid inks. Aqueous inks remain in a liquid form when stored prior to being used in imaging operations. Solid ink or phase change inks typically have a solid form at room temperature, either as pellets or as ink sticks of colored ink, which are inserted into feed channels in a printer through openings to the channels. After the ink sticks are fed into the printer, they are urged by gravity or a mechanical actuator to a heater assembly of the printer. The heater assembly includes a heater and a melt plate. The heater, which converts electrical energy into heat, is attached to the melt plate to heat the melt plate to a temperature that melts an ink stick coming into contact with the melt plate. The melt plate can be oriented to drip melted ink into a reservoir and the ink stored in the reservoir continues to be heated while awaiting subsequent use.\nThe reservoir can be attached directly to the printhead or can be a separate subsystem. With either approach, fluid couplings in the printer supply the liquid ink held in each reservoir of colored ink to one or more printheads in the inkjet printing apparatus. The liquid ink is pumped from the reservoir to a manifold in the inkjet printing apparatus. The manifolds are connected through channels to numerous individual inkjet ejectors. As the inkjets in the printheads eject ink onto a receiving medium or imaging member, the action of the diaphragms in the inkjets pulls ink from the manifold. Various embodiments of inkjets include piezoelectric and thermal devices that are selectively activated by a controller with an electrical firing signal.\nPhase change ink printers often include one or more heaters that maintain a supply of phase change ink in a liquid state for use during printing operations. Typically, the heaters are electrical heaters that consume electrical energy to maintain the phase change ink in a liquid phase. In order to reduce energy usage, phase change ink printers deactivate various components, including heaters, in the printer during a sleep mode to conserve energy. The ink held in the printheads and inkjets cools and solidifies in some sleep modes.\nWhile sleep modes enable a printer to operate with reduced electrical energy consumption, the solidification of phase change ink within the printer presents difficulties to printing high quality documents when the printer emerges from sleep mode. As phase change ink within an inkjet printing apparatus cools and solidifies, the ink contracts and air enters the pressure chambers and fluid conduits within the printheads. As the solidified ink heats and liquefies during a subsequent warm-up process, the air forms bubbles in the liquefied ink that can prevent inkjets in the printheads from operating reliably. Additionally, during the warm-up process, both the ink and air bubbles expand due to the heat applied to the printheads. The expanding air bubbles may force some ink through the ejector nozzles, which is referred to as “drooling.” The drooled ink can contaminate other nozzles in the printheads or separate from the printheads and produce errant marks on the image receiving member.\nTo eliminate air bubbles in the liquefied ink within the printheads and to clear contaminants from the inkjet nozzles and external face of each printhead, the inkjet printing apparatus undergoes a “purge” operation where pressure applied to the printheads urges the liquid ink and the air bubbles through the nozzles of the inkjets. In a purge operation, the inkjets emit a stream of ink that flows down the face of the printhead and is collected in a waste ink receptacle instead of being ejected as individual ink drops. The purge operation removes air bubbles from the inkjets in the printheads and other fluid conduits in the inkjet printing apparatus.\nIn some printing apparatus designs, a wiping operation occurs after the purge operation. In a wiping operation, a wiper blade engages the face of a printhead and moves across the printhead face, including the inkjet nozzles. The wiper blade cleans residual ink and contaminants on the face of the printhead from the purge operation. The wiping operation maintains the meniscus formed between the liquid ink and nozzle in each of the inkjets in the printhead. The meniscus may be broken if the liquid ink contacts a contaminant or another mass of liquefied ink on the face of the printhead. The wiping operation clears the contaminants to enable each inkjet to maintain the meniscus for reliable operation.\nIn existing printers, the purged ink and ink from a wiping operation is typically collected in a waste reservoir and is eventually discarded. In printers that enter sleep modes more often to reduce electrical energy consumption, the number of purge cycles and the corresponding amount of discarded ink increases. A desirable improvement to phase change ink-jet printers would reduce or eliminate discarded ink produced during purge cycles."}
{"text": "Remarkable growth in the demand for communications products and services, and especially in the requirements for portable communications devices, has driven consumer requirements for low-cost, small-form-factor, low-power RF (radio frequency) transceivers. In addition, the development of state-of-the-art wireless applications has encouraged consumers to expect both the convenience of extended connectivity and the benefit of enhanced services. RF transceivers that operate in compliance with multiple prevailing standards are instrumental, if not required, in the satisfaction of these objectives. In this regard, the capabilities of CMOS (complementary metal/oxide/semiconductor) and BiCMOS (bipolar/CMOS) VLSI (very large scale integration) technology are particularly well suited to the accommodation of very aggressive levels of mixed-signal integration, as well as to the provision of increasing functionality in a single-chip RF integrated circuit (IC) device.\nHowever, the increasing density to which integrated RF circuit blocks are packaged in the same, or neighboring, integrated circuit devices has been attended by a number of operational challenges. Prominent among them is an aggravated susceptibility to electromagnetic interference (EMI) that may be propagated between circuits and devices. For example, with respect to densely packaged integrated circuit devices, and with respect to RF integrated circuits in particular, a current circulating in one loop in a device, or on a circuit board, may induce an interfering voltage at the input of another circuit or device. The likelihood of interference is exacerbated when the circulating current is an RF current that propagates to a circuit that is required to accept signals at the frequency of the EMI.\nConventionally, an interfering voltage may be reduced by one or more of the following methods: minimizing the loop area of either the transmitting or the receiving loop, increasing the distance between the transmitting and receiving loops, or shielding either or both the transmitting or the receiving loop. In the situation where an integrated circuit may be part of the receiving loop, it may be impracticable either to further minimize loop area or to provide shielding as a mechanism to reduce the induced interference. That is, the input loop of an integrated circuit consists essentially of a lead frame, having substantially fixed dimensions, and bond wires connecting the lead frame to the integrated circuit die. The dimensions and overall area of this part of the input loop are tightly controlled and fixed, but may not be reduced to zero. Furthermore, standard IC packaging does not provide even minimally effective shielding for this part of the input loop.\nAccordingly, what is desired is a technique to improve the immunity of a receiving (IC) device from interfering signals that may arise in the proximity of the device."}
{"text": "As computer and other electrical equipment continue to drop in price, the manufacturers of storage devices, such as memory devices and hard drives, are forced to lower the cost of their components. At the same time, computer, video game, television and other electrical device markets are requiring increasingly larger amounts of memory to store images, photographs, videos, movies, music and other storage intensive data. Thus, besides reducing cost, manufacturers of storage devices must also increase the storage density of their devices. This trend of increasing memory storage density while reducing cost required to create the storage has been on-going for many years, and even optical storage such as CD-ROM, CD-R, CD-R/W, DVD, and DVD-R variants are being challenged by device size limitations and cost. There is accordingly a need for economical, high capacity memory structures."}
{"text": "The unbroken popularity of the World Wide Web and its annual increase in size requires increasingly larger switching fabrics. To meet today's switching requirement of 1 Tb/s throughput and above, switches are growing not only in terms of port speed, but also in number of ports. With 5 Gb/s links, port sizes of 64, and today's link and fabrication technologies, multichip and even multishelf solutions are necessary. As a consequence, the input adapters with virtual output queuing (VOQ) are at a distance from the switching device such that the roundtrip time (RTT) bandwidth product is significantly greater than 1. On the other hand switch packet sizes in terms of bytes remain constant. This means that there are more packets on the link, which requires larger memories in the switch, for it to be able to accept all packets in transit. This is a necessity in order that the switch be work-conserving or loss-less, depending on the link-level flow-control scheme used.\nA large amount of the expenses invested in switch design goes into link technology, i.e. links, cables, and connectors. Therefore bandwidth is expensive, but however the bandwidth is not used efficiently at present.\nIBM's PRIZMA switch-chip family uses a grant flow-control scheme that returns to each input adapter a vector of stop and go signals per output. This scheme is described in the publication “A Combined Input- and Output-Queued Packet-Switch System Based on PRIZMA Switch-on-a-Chip Technology” by C. Minkenberg and T. Engbersen in IEEE Commun. Mag., vol. 38, no. 12, December 2000, pp. 70-77. For fabric sizes of N=64 for instance, with N the number of switch ports, a vector length and hence a flow control bandwidth of 8 bytes per packet cycle would be required. If there are not potentially complex optimization schemes applied, the grant flow-control scheme will prevent switches to grow to larger sizes.\nThe Atlas switch as described by G. Kornaros et al. in the publication “Implementation of ATLAS I: a Single-Chip ATM Switch with Backpressure” in Proc. IEEE Hot Interconnects VI Symposium, Stanford, Calif., USA, 13-15, Aug. 1998, pp. 85-96 represents the category of switches that use a flow-control scheme based on credits. The flow-control bandwidth is designed to return two credits per packet cycle. The storage, serialization, and return of credits is performed per input using a so-called credit-out FIFO (first-in first-out memory). The FIFO must be large enough to hold all credits that are allowed in the worst case to circulate per adapter/switch input pair. For reasons of correctness the FIFO size scales with the number of ports and the memory size assigned per switch input/output pair. Therefore, the FIFO size roughly scales with O(MN), where M signifies the assigned memory size per memory point, and N the number of switch ports. More importantly, each FIFO must provide N write accesses per packet cycle, because each switch row can have up to N parallel departures. For N≧16 this presents tough hardware design challenges. For ever larger switches it is nearly an impossible task. Therefore the FIFO concept for storage and serialization is not a scalable solution.\nFurther problems arise from the limited bandwidth of the in-band flow-control channel. An out-of-band flow-control is prohibitively expensive for terabit and betabit solutions. Flow-control bandwidth becomes a real bottleneck if one scales existing switches to ever larger sizes realized as single or multistage fabrics. Furthermore, for scalable multistage fabrics, there are severe restrictions that enforce a number of flow-control events per packet-cycle in the channel.\nFrom the above follows that there is still a need in the art for a new kind of flow-control mechanism which is performant, efficient, robust, scalable, and has the potential to be used in future switching fabrics. The mechanism should be suitable for various switch environments such as communication systems and multiprocessor interconnects. It is therefore an object of the present invention to provide an improved flow-control mechanism for high and efficient packet throughput."}
{"text": "For a surface layer used for a display and the like, (1) those obtained by a method generally referred to as a dry method, i.e., making a dielectric multilayer film in a vapor phase process and realizing a low reflectance with an optical interference effect and (2) those obtained by a method generally referred to as a wet method, i.e., coating a low refractive index material on a substrate film have been used. As a technology which is quite different principally from (1) and (2), (3) it has been known that the low reflectance can be realized by providing a fine structure to the surface (Patent Documents 1 to 10).\nIn general, such a surface layer has required not only an antireflection performance of light and an improved performance of the light permeability, but also a certain mechanical strength to withstand abrasion and scratches in practical use.\nHowever, for the surface layer having the surface fine structure described in the above (3), the good antireflection performance is obtained while the mechanical strength such as surface scratch resistance is insufficient. Thus, the surface layer easily abrades away and is easily scratched. Therefore, it has not come in practical use yet.\nFor example, in Patent Documents 1 to 10, materials for such an antireflection film are listed, and a (meth)acrylate compound is described therein to be used as a polymerizable resin. However, the materials listed there are quite usual materials for forming ordinary polymer films, and they are not the materials which are discussed from the aspect of the present invention. Namely, Patent Documents 1 to 10 do not describe the surface layer having the special surface fine structure described in (3), which is made practical for the mechanical strength such as surface scratch resistance by selecting these materials.\nFurthermore, the antireflection film having the fine structure on the surface in the above described (3) has the special fine structure on the surface so as to suitably prevent the reflection. Thus, for the materials to be selected for the antireflection film, a special physical property is required. However, what kind of physical property is required has been scarcely known.\nPatent Document 1: JP Sho-50-070040-A\nPatent Document 2: JP Hei-9-193332-A\nPatent Document 3: JP 2003-162205-A\nPatent Document 4: JP 2003-215314-A\nPatent Document 5: JP 2003-240903-A\nPatent Document 6: 2004-004515-A\nPatent Document 7: 2004-059820-A\nPatent Document 8: 2004-059822-A\nPatent Document 9: 2005-010231-A\nPatent Document 10: 2005-092099-A"}
{"text": "1. Field of the Disclosure\nThe subject disclosure relates generally to oilfield drilling, and more particularly to bottom hole assemblies and tools for orienting a bottom hole assembly (BHA).\n2. Background of the Related Art\nIn conventional drilling, the BHA is lowered into the wellbore using jointed drill pipes or coiled tubing. Often the BHA includes a mud motor, directional drilling and measuring equipment, measurements-while-drilling tools, logging-while-drilling tools and other specialized devices. A simple BHA having a drill bit, various crossovers, and drill collars is relatively inexpensive, costing a few hundred thousand US dollars, while a complex BHA costs ten times or more than that amount.\nMany drilling operations require directional control so as to position the well along a particular trajectory into a formation. Directional control, also referred to as “directional drilling,” is accomplished using special BHA configurations, instruments to measure the path of the wellbore in three-dimensional space, data links to communicate measurements taken downhole to the surface, mud motors, and special BHA components and drill bits. The directional driller can use drilling parameters such as weight-on-bit and rotary speed to deflect the bit away from the axis of the existing wellbore. In some cases, e.g. when drilling into steeply dipping formations or when experiencing an unpredictable deviation in conventional drilling operations, directional-drilling techniques may be employed to ensure that the hole is drilled vertically.\nDirection control is most commonly accomplished through the use of a bend near the bit in a downhole steerable mud motor. The bend points the bit in a direction different from the axis of the wellbore when the entire drill string is not rotating. By pumping mud through the mud motor the bit rotates (though the drill string itself does not), allowing the bit alone to drill in the direction to which it points. When a particular wellbore direction is achieved, the new direction may be maintained by then rotating the entire drill string, including the bent section, so that the drill bit does not drill in a direction away from the intended wellbore axis, but instead sweeps around, bringing its direction in line with the existing wellbore. As it is well known by those skilled in the art, a drill bit has a tendency to stray from its intended drilling direction, a phenomenon known as “drill bit walk”. A device for addressing drill bit walk is shown in U.S. Pat. No. 7,610,970 to Sihler et al. issued Nov. 3, 2009, which is incorporated herein by reference.\nThe use of coiled tubing with downhole mud motors to turn the drill bit to deepen a wellbore is another form of drilling, one which proceeds quickly compared to using a jointed pipe drilling rig. By using coiled tubing, the connection time required with rotary drilling is eliminated. Coiled tube drilling is economical in several applications, such as drilling narrow wells, working in areas where a small rig footprint is essential, or when reentering wells for work-over operations.\nIn coiled tubing drilling, a BHA with a mud motor is attached to the end of a coiled tubing string. Typically, the mud motor has a fixed or adjustable bend housing to drill deviated holes. Because the coiled tubing is unable to rotate from surface, a so called orienter tool is used as part of the BHA to “orient” the bend of the mud motor into the desired direction. There exists a multitude of different designs for the drive systems of such tools. Some designs support continuous rotation such as electric motor and gearbox drives, while others only permit rotation by a certain limited angle."}
{"text": "1. Technical Field\nThe present invention relates to supply of a liquid.\n2. Related Art\nDevices which supply liquids to handpieces intermittently ejecting the liquids are known (see JP-A-2014-95353). The liquid supply device disclosed in JP-A-2014-95353 includes two plunger pumps. Pulsation of a liquid to be supplied is reduced by causing the two plunger pumps to alternately supply the liquid.\nIn the case of the technology of the foregoing related art, it takes some labor and time to perform maintenance of the plunger pumps. In particular, when the plunger pumps are used for medical treatment, problems particularly occur in cases in which exchange or sterilization is necessary."}
{"text": "Recent advances in silicone technology have resulted in a variety of valve designs for children's sippy cups. Although many existing designs are effective in preventing the escape of liquid from a container without the application of negative pressure from a user, these valves are generally permanently installed in the cup, thereby reducing sanitation and longevity of the cup. Particularly when used by children or the infirmed, sterilization of all components of a cup is extremely important, but cannot be accomplished without each component's removal.\nPermanent valve installations in such cups can cause mold and bacteria buildup, which can eventually pose significant risks to the user. Accordingly, cups having permanent valve installations may be used a limited number of times before customary cleaning techniques are no longer able to thoroughly sanitize the valve components, requiring the user to dispose of the cup. Such requisite disposal has obvious economic and environmental implications to the users and consumers generally, in addition to the significant distress disposal of a child's favorite cup can impose on familial relations.\nWhen sippy cup designs provide for a removable valve assembly, however, the assembly is generally constructed of a single material that, while flexible, is generally easily torn with repeated extractions and insertions of the valve assembly. Given the repeated washings required for utensils and containers used by children, in particular, durability is a significant aspect for consumers of such products. Accordingly, existing valve assemblies fail to provide the removability—desirable from a sanitation perspective—and durability—preferable from an economic and environmental perspective—deemed necessary to modern consumers, and possible to manufacturers given new advances in overmolding processes.\nThere is disclosed herein an improved sippy cup assembly which avoids the disadvantages of prior devices while affording additional structural and operating advantages."}
{"text": "This invention relates to a cold-air generating device, and in particular, to a cold-air generating device which employs a plurality of evaporators and compressors, thereby improving the evaporator heat exchange efficiency and making it possible to reduce the level of starting current needed.\nAn example of conventional cold-air generating device of this type is disclosed in Japanese Utility-Model Publication No. 48-22155.\nThe conventional device disclosed therein is equipped with a refrigerant compressor, a condenser, a liquid tank, an evaporator, a reheater, a capillary tube and an electromagnetic valve for changing cycles. In its refrigeration cycle, which is designed to enable switching between dehumidification and air-conditioning operations, a by-pass circuit is formed in parallel with the liquid tank so that, during dehumidification, the refrigerant will be conveyed to the reheater instead of being passed into the liquid tank.\nIt has been the practice with conventional cold-air generating devices of the above-described type to use a blowing compressor when the air is to be refrigerated down to a temperature of 0.degree. C. or below (-20.degree. C. to -80.degree. C.).\nSuch a compressor, however, involves noise and vibration, and requires additional capacity, which makes the cold-air generating device unsuited for medical use.\nIn regard to power capacity, it may be noticed that, given a limited power source of AC100 V, the output of commercially available blowing compressors is 0.75 KW or less.\nIf a blowing compressor with an output of 0.75 KW is used, the starting current required will be 50 A or so. Consequently, a dedicated power source will have to be provided in order to start this blowing compressor. This requires extra electrical work.\nIf such extra electrical work is not to be performed, the output of the blowing compressor must be below 0.4 KW. In the case of a commercially available blowing compressor with an output of 0.4 KW, the airflow volume is a little less than 50 l/min, a value which is too low for use in a medical apparatus for cryotherapy.\nTo refrigerate air down to a temperature of 0.degree. C. or below (-20.degree. C. to -80.degree. C.), a single evaporator is not enough because almost all the water contained in the air will then turn into ice and adhere to the surface of the evaporator. If air at a temperature of 20.degree. C. and with a relative humidity of 60% is cooled down to -30.degree. C., about 10 g of ice per 1 m.sup.3 air adheres to the evaporator surface. If the air is allowed to flow at a rate of 100 l/min, 60 g of ice per hour will adhere to the surface. Accordingly, if the air is allowed to flow for a long time, the smoothness of the flow will be impaired, resulting in the evaporator providing a poor heat exchange efficiency.\nBesides this, provision of a single evaporator requires that its size be considerably large; the refrigerating compressor should also be large. Since only one refrigerating compressor is provided, the current required for starting it is relatively large as compared with the case where two refrigerating compressors are provided (the starting current that flows is about four or five times as large as the rated current), and poor starting may be the result if the power capacity is relatively small. If starting problems occur, some extra electrical work has to be conducted for the purpose of raising the power capacity."}
{"text": "A device and a method of this type with scattering correction are known from U.S. Pat. No. 5,666,391 A. The known device and the known method allow x-ray scatter to be corrected within the framework of computer tomography.\nComputer Tomography (=CT) is a well-known method with which the three-dimensional reconstruction of the tissue density distribution, or to put it more precisely of the linear coefficient of attenuation of x-rays, of the inside of the body is possible with high spatial resolution and great quantitative accuracy of density resolution. The spatial resolution is typically <=1 mm and the density resolution typically amounts to a few Hounsfield Units (1 HU= 1/1000 of the density of water).\nWith two-dimensional detectors with a typical size of around 40 cm×40 cm such as for example x-ray image amplifiers or Flat-Panel Detectors (=FPD), It is possible, with just a single half rotation, which extends over more than 180 degrees and the cone angle, to record and to reconstruct an image of a large volume area, (=Conebeam CT=CBCT)."}
{"text": "The discussion below relates generally to the field of enterprise resource planning (ERP) software. ERP software is business management software that allows an organization to use a system of integrated applications to manage business. ERP software can integrate all facets of an operation, including development, manufacturing, sales and marketing. ERP software can include many enterprise software modules that are individually purchased, based on what best meets the specific needs and technical capabilities of the organization. Each ERP module can be focused on one area of business processes, such as product development or marketing. Some of the more common ERP modules can include those for product planning, material purchasing, inventory control, distribution, accounting, marketing, finance and human resources (HR). As the ERP methodology has become more popular, software applications have emerged to help business managers implement ERP in other business activities. Applications may incorporate modules for customer relationship management (CRM) and business intelligence. These modules can be presented as a single unified package. A central repository can store all information that is shared by all the various ERP facets in order to smooth the flow of data across the organization.\nOracle's PeopleSoft supplier relationship management is an integrated suite of procurement applications that can reduce supply management costs. PeopleSoft supplier relationship management can reduce spending on goods and services, streamline procure-to-pay processes, and drive policy compliance. PeopleSoft supplier relationship management provides packaged integration of procurement functions with financial management, human capital management, and ERP suites. PeopleSoft supplier relationship management can give customers the flexibility to leverage applications on-demand, on-premises, or in any combination needed to achieve procurement objectives. PeopleSoft supplier relationship management can integrate procurement with HR in order to provide total workforce management of both contingent and full-time workers.\nOrganizations recognize the importance of strategic procurement practices as a key component of improved profitability. Request-for-quote and request-for-proposal (collectively RFx) activities can focus on finding the best value and lowest total cost supplier rather than just searching for the lowest price. Tools such as online auctions can be used to increase supplier competition and to benefit to organizations across a growing number of industries.\nStrategic sourcing refers to the process of determining the best suppliers for needed goods or services, and the conditions under which business is to be awarded to those suppliers. The strategic element of sourcing can come from the ability to apply different methodologies to create the most value in an organization's supply chain. This value can come in the form of pricing, warranty terms, quality, delivery accuracy, or other factors. Strategic sourcing enables an organization to quantify and balance such requirements so that the organization can achieve its specific business objectives.\nUsing sourcing tools and strategies, organizations can make buying decisions as part of an overall strategy for achieving business goals, with a view toward building long-term relationships with key suppliers. In an era of ever-increasing demands for cost control and higher operational performance, procurement has become a mission-critical operation. A complete solution built on interact-based portal technologies, reverse auction capabilities, and proposal analysis tools can deliver significant value to an organization.\nPeopleSoft Strategic Sourcing can enable an organization to efficiently onboard bidders with self-service registration, reduce procurement costs through competitive bidding and reverse auctions, align its purchasing objectives and execution with overall business strategy, automate and control the RFx and auction process both inside and outside the enterprise, leverage existing supplier and customer relationships while reaching out to new trading partners, track bidder participation and results effectively, retain its knowledge and strategies in an application for future use and performance analysis, perform award analysis using optimization to help determine optimal awards, plan sourcing activities using planning tools that help track progress, and estimate the total cost of transacting with suppliers beyond just price.\nSupplier contract management has become even more overwhelming in the face of new business challenges such as globalization, outsourcing, and increased regulatory pressure. Today, most companies record some financial details about their contracts in their purchasing system, but many still store the actual contract document offline (for example, in a file cabinet) with no connection to the execution and compliance activity. Authoring remains a manual process lacking many if any controls. Terms negotiations can be cumbersome and time-consuming, and incomplete adoption and lackluster enforcement can prevent even the best-negotiated contracts from realizing expected savings."}
{"text": "1. Field of the Invention\nThis invention relates to true phase matched parametric generation of light in monolithically integrated intersubband (ISB) semiconductor optical devices in general and, more particularly, to true phase matched second harmonic (SH) generation of light in monolithically integrated quantum cascade (QC) semiconductor optical devices.\n2. Discussion of the Related Art\nBy the phrase parametric generation we mean any optical process by which light of one center frequency is made incident on or is generated inside a nonlinear optical region (e.g., a suitably designed quantum well structure; a bulk Group III-V compound semiconductor; LiNbO3) to generate light of a different center frequency and often at harmonics thereof. As such, parametric generation includes two-wave interactions such as second harmonic (SH) generation and three-wave interactions such as sum and difference frequency generation, as well as interactions involving more than three waves (e.g., four-wave mixing).\nIn the prior art realm of non-integrated nonlinear optical devices, it is known to utilize discrete optical pump sources and discrete nonlinear bodies to generate light at a frequency different from that of the pump source. Thus, for example, in IEEE Photonics Technology Lett., Vol. 12, No. 5, pp. 486–488 (2000) Chowdhury et al. analyze continuously phase matched M-waveguides in LiNbO3 for second-order nonlinear up-conversion. Implicit in the analysis is that the pump source is external to and not integrated with the LiNbO3 waveguide. Similarly, in Appl. Phys. Lett., Vol. 83, No. 4, pp. 620–622 (2003) Moutzouris et al. describe SH generation through optimized modal phase matching in GaAs/AlGaAs waveguides. Explicit in the experimental work is the use of a 1.55 μm optical pump source external to and not integrated with the semiconductor waveguide. Inherently, however, these types of discrete parametric devices suffer from relatively low nonlinear conversion efficiency as well as optical coupling losses between the source and nonlinear body.\nIn contrast, bodies of nonlinear optical semiconductor materials show promise for achieving high-efficiency, nonlinear conversion of light when monolithically integrated within compact, injection-pumped semiconductor light emitters. These integrated optical devices would extend the spectral range accessible to semiconductor optical sources and could find applications in fields ranging from high-resolution chemical sensing to quantum cryptography.\nAlthough resonant ISB transitions have been studied extensively as a source of nonlinear (or parametric) light generation, so far practical applications have been limited by two obstacles: the lack of both a powerful integrated pump source and a tunable phase matching scheme. [See, for example, Gurnick et al., IEEE J. Quantum Electron., Vol. QE-19, No. 5, pp. 791–794 (1983); Rosencher et al., Science, Vol. 271, pp. 168–273 (1996); Capasso et al., IEEE J. Quantum Electron., Vol. 30, No. 5, pp. 1313–1326 (1994); and Vurgaftman et al., IEEE J. Quantum Electron., Vol. 32, No. 8, pp. 1334–1346 (1996), all of which are incorporated herein by reference.] Only recently has the first obstacle been overcome by monolithically integrating the nonlinear ISB transitions within the structure of a QC laser. [See, for example, Owschimikow et al., Phys. Rev. Lett., Vol. 90, No. 4, pp. 043902-(1–4) (2003) and Gmachl et al., IEEE J. Quantum Electron., Vol. 39, No. 11, pp. 1345–1355 (2003), both of which are incorporated herein by reference.] On the other hand, the challenge of phase-matching the modes of the pump (laser) light and the modes of the parametric (e.g., SH) light continues to be an obstacle to achieving high efficiency parametric conversion.\nCommon techniques for phase matching, such as birefringence phase-matching or quasi phase matching, are not practical for ISB lasers. [Regarding the former, see, for example, Fiore et al., Nature, Vol. 391, pp. 463–466 (1998), which is incorporated herein by reference.] Moreover, the various schemes proposed for parametric (e.g., SH) phase matching in asymmetric double quantum well structures cannot be easily applied to ISB lasers because of the intrinsic waveguide dispersion, weak voltage tunability, and strict current transport requirements. [These schemes are described by Vurgaftman et al., supra, Meyer et al., Appl. Phys. Lett., Vol. 67, No. 5, pp. 608–610 (1995) and Vodopyanov et al., Appl. Phys. Lett., Vol. 72, No. 21, pp. 2654–2656 (1998), all of which are incorporated herein by reference.]\nThus, a need remains in the art for a monolithically integrated semiconductor optical source that exhibits true phase matching and efficient parametric light generation.\nBy the phrase true phase matching we mean that the nonlinear body is characterized in that the effective refractive indices of the mode of the pump light and that of the parametrically generated light are essentially equal to one another. True phase matching is to be distinguished from prior art techniques such as quasi phase matching, which relies on periodically resetting the phase of the parametric light at the layer interfaces of a multilayer nonlinear body, and birefringence phase matching, which utilizes the different indices of refraction of ordinary and extraordinary waves to make the effective refractive index of the pump light equal to that of the parametric light."}
{"text": "This disclosure relates to powertrain systems and, more particularly, to cylinder deactivation powertrain systems.\nConventionally, applying cylinder deactivation on piston engine powertrains provides improved fuel economy over a combination of speeds and loads limited by excessive torque pulses. However, the limited speed-load range of those piston engine powertrains, where cylinder deactivation can be used, reduce the overall improvement in fuel economy to such a degree that few production vehicles employ cylinder deactivation. When cylinder deactivation is employed, one or more cylinders are fired and then skipped, or deactivated, and a time lapse occurs until the next cylinder fires. The resulting time lapse translates into roughness that can become a customer dissatisfier. This irregular or staggered torque output reduces the fuel economy efficiency otherwise available from cylinder deactivation. In addition, the exhaust emissions generally increase due to the inconsistent fuel delivery and combustion.\nAccordingly, there exists a need to selectively fire cylinders without causing objectionable torque pulses and increasing emissions.\nThe drawbacks and disadvantages of the prior art are overcome by the embodiments of the powertrain system described herein. The powertrain system comprises an engine coupled to a drive shaft, and a torque-smoothing device that is configured to sometimes assist the drive shaft by adding energy, and at other times to move energy from the drive shaft to the energy storage device. The engine is also equipped with a cylinder deactivation system configured to deactivate one or more engine cylinders, and a direct cylinder injection system configured to deliver fuel to firing cylinders. An engine management system is configured to provide control, diagnostic and maintenance operations for the powertrain system.\nThe method for operating the internal combustion engine comprises activating the engine, and deactivating one or more engine cylinders. When the engine generates an inconsistent torque output, a torque-smoothing system is activated to smooth the torque output of the engine."}
{"text": "The present disclosures relate generally to semiconductor memories, and more particularly, to semiconductor memories with recessed devices.\nAs bitcell size for static random access memories (SRAMs) continues to scale to smaller sizes, the bitcell current (Icell) performance degrades. In addition, static noise margin (SNM) variation of the bitcell increases. Together, the degradation of Icell performance and the increase in SNM variation limits the low supply voltage (Vdd) operation of the bitcell.\nSome SRAM devices have been known to achieve a higher bitcell current, however, the SRAM's beta ratio (i.e., the ratio of strength of the pull-down device to the pass device) suffers negatively. As a result, the adverse impact on the SRAM beta ratio degrades the SNM and makes the bitcell unstable at low Vdd operation. Still further, the bitcell may be unstable even at nominal Vdd operation. Furthermore, with scaling to smaller and smaller sizes, transistor threshold voltage (Vt) variation increases due to dopant fluctuations and variations in gate length.\nAccordingly, it would be desirable to provide an improved memory for overcoming the problems in the art, as discussed above."}
{"text": "The invention concerns a device and a method to replace at least one component of a pourable medium. Furthermore, the invention concerns the use of such a device in a printer or copier.\nExamples of systems in which pourable medium are used are electro-photographic or ionographic printers or copiers. The pourable medium used therein is a developer, especially a two-component developer that contains a toner and a magnetic carrier respectively in particle form.\nIn printers or copiers operating according to the electro-photographic principle, charge images of characters (such as texts or drawings) to be printed are generated on an intermediate carrier, for example a drum with an electro-photographic layer. The intermediate carrier with the charge images is directed into a development station in which the charge images are developed and thereby made visible. Upon the development, the toner particles of a two-component developer are conveyed in the direction of the intermediate carrier and are electrically attracted onto the intermediate carrier by the charge carrier. The carrier particles are conveyed back to the development station. The developed charge images are transferred and fixed upon a recording medium, for example a paper web. The toner material within the two-component developer is used in the development valve. Fresh toner material is refilled corresponding to this use. The carrier material conveyed back is subject to an aging process, whereby the print quality or copy quality can suffer. Therefore, it is also necessary to renew the carrier material little by little.\nA printer device or copier device is specified in WO 98/39691 (of which reference is made here to its entire disclosure) in which a recording medium can be printed upon in monochrome and/or color. A plurality of developer stations are thereby provided for which, respectively, one color is associated for one colored print. In order to maintain the print quality in printers of such a type, a replacement must occur at specific time intervals due to the expiration of the developer or the expired components contained therein in the form of the carrier. For example, this could hitherto occur such that development stations are completely emptied and subsequently filled, which is very time-consuming and is associated with an interruption of the operation of the printer or copier.\nReplacement of the developer with expired carrier with fresh developer can also occur such that developer stations are supplied with fresh toner and fresh carrier separate from one another, such that a definite amount of the toner and carrier ensuring the print quality or copy quality is always present in the developer, such that the printer operation or copier operation must not be interrupted.\nA device for the replacement of expired developer with fresh developer in a development station of an electro-photographic device is known from U.S. Pat. No. 4,358,196. This known device comprises two metering devices, via one of which fresh developer is supplied to the developer station, and via the other of which expired developer is removed from the development station, such that a developer ensuring good print quality is always present in the development station.\nGiven replacement of the developer with expired carrier, it is important that a specific volume of developer with a defined toner concentration and carrier concentration (meaning a defined proportion of toner to carrier) is always present in a development station. Over the entire lifespan of a printing device, a total fill quantity of developer guaranteeing proper functioning must always be present. The toner concentration can be regulated in its own regulation cycle, such that an amount of toner ensuring the print quality is always present in the development station. The regulation can occur such that the toner concentration is detected and regulated either directly by means of a sensor, or indirectly via detection of the inking of the printing. However, given direct regulation, the toner concentration cannot be precisely determined due to inevitable control deviations, and given indirect regulation, the toner concentration can in part not be determined at all, whereby losses in quality can arise.\nTypically, the ratio of toner to carrier for the volume of the supplied two-component developer is known. However, the ratio of toner to carrier in the volume of the removed developer can vary from the ratio in the supplied developer, because the toner concentration in the development station can vary higher or lower than the toner concentration of the supplied developer due to the control deviations. However, because both volumes are equal, the proportion of the carrier in the removed developer can vary higher or lower than the proportion of the carrier in the supplied developer.\nIf the replacement of the developer is implemented uniformly and constantly, as is the case in the device according to the US patent specification cited above, the carrier quantity in the developer station can drift ever further higher or lower, given a large number of alternating cycles. This means that the development station can overflow or run empty in the extreme case, whereby the printer or copier is incapable of functioning.\nA device and a method for the replacement of at least one component of a pourable medium is known from EP 0 043 505 A2. The pourable medium with the expired component is removed from a container, and a quantity of fresh pourable medium is supplied into the container. Two metering devices are provided, from the first of which the developer station is constantly supplied fresh developer, and from the other of which used developer is constantly removed from the developer station in an equal quantity.\nA developer station is known from U.S. Pat. No. 5,450,178, in which a mix of toner and carrier particles is present in an upper chamber according to the exemplary embodiment according to FIG. 4. The proportion of toner is higher than is necessary for the typical usage. The mix from the upper chamber is transferred little by little into a container, whereby an expired mixture in this container is dispensed into a lower chamber, dependent on the signal of a sensor that analyzes the contents in the container. In this manner, the total quantity of the mixture can be maintained at a constant value.\nA developer unit is known from U.S. 5,095,338, in which a mixture of toner particles and carrier particles is present in a chamber (compare in particular FIG. 2). A quantity of carrier particles and toner particles is supplied from the container to the development material in the chamber. The ratio of toner particles to carrier particles in the container is greater than in the chamber, such that expired developer material is continually renewed. A sensor system observes the composition of the development material in the chamber and initiates via a controller the supply of fresh developer material.\nA copier with at least two developer stations for at least two different colors is known from EP 0 043 505 A2, in which a fresh developer component is supplied from a common reservoir container to these developer stations at the respective supply intake. This developer is a liquid developer.\nReference is made to the documents JP 59 100 471 A, JP 2000 122 424 A, and WO 01/06325 A as further relevant prior art."}
{"text": "1. Field of the Invention\nThe present invention relates to an infrared achromatic waveplate capable of improved response as a function of wavelength in a simple, small and inexpensive configuration of waveplates.\n2. Discussion of Background\nMany optical systems including those for spectropolarimetry, laser polarimetry, laser spectroscopy, and ellipsometry have a need for converting light between polarization states and a need to analyze polarized light. The design of polarimeters requires polarization elements whose properties satisfy a number of criteria including the very important criteria that the polarization properties need to be substantially constant over a range of wavelength of interest. Additional constraints in this area include a reasonable element size and proper positioning of the light beam exiting from the element and of course the cost of the element.\nLiquids and amorphous solids such as glass and crystalline solids have a cubic symmetry which normally show a behavior whereby the speed of light and the index of refraction is independent of the direction of propagation in the medium and is independent of the state of the polarization of the light. These types of elements are said to be optically isotropic. Other crystalline solids, which induce birefringent behavior, are optically anisotropic. Of course, solids may be anisotropic in many of their properties, but it is the optical anisotropy of a material which is used in order to provide the \"double refraction\" of a beam. The two emerging beams from an optically anisotropic material are plane-polarized beams with their planes of vibration at right angles to each other.\nThe conversion of light between polarization states and the analysis of polarized light has traditionally involved the use of birefringent materials wherein a light beam incident on a birefringent material is divided into two orthogonal polarization components. A retarder can then shift the phase or in other words retard the phase of one of these two orthogonal polarization components with respect to the other component. In a birefringent material, the index of refraction depends on the polarization state of the light beam.\nThe most appropriate way that these anisotropic or birefringent materials are used involves the exploiting of the dependency of the index of refraction on the polarization state of the incoming light beam. A phase shift is introduced between the polarization state aligned with the fast axis of the birefringent material, where the index of refraction is the smallest and the polarization state aligned along the slow axis of the material, where the index of refraction is at its highest value.\nWhen plane-polarized light falls at normal incidence on a slab or piece of anisotropic material so that the optic axis is parallel to the face of the slab, the two waves which emerge are plane-polarized at right angles to each other. Because the waves travel through the material at different speeds, there will be a phase difference between the two waves when they emerge from the material. If the material thickness is chosen so that for a particular frequency of light the phase angle or \"phase change between the two waves\" is 90.degree., the slab or piece of material is called a quarter-waveplate. If linearly polarized light is incident on this quarter waveplate with its plane of polarization oriented at .+-.45.degree. to the fast axis, the emerging light is said to be circularly polarized.\nTraditionally then, the proper thickness of the material was chosen in order to obtain the desired retardance. However, these prior art designs are very sensitive to small changes in wavelength of the incident beam and thus are not suitable for many optical systems where broadband light is used.\nWhile other designs have better response as a function of wavelength, they involve complicated, large and expensive devices as for example in a design utilizing a modified Fresnel rhomb which is 4 inches long which of course exceeds the requirements for size.\nStill other designs cause the exit beam to be shifted from the path of the incident beam which also makes these designs not appropriate for such polarimeter usage.\nThus, there is a specific need for an achromatic infrared retarder in which the polarization properties of the element is substantially constant over a particular range of wavelength and wherein the elements are small in size and produce a light beam which is properly positioned upon emergence from the element."}
{"text": "The parent application, Ser. No. 08/155,252, which is incorporated herein by reference, teaches a fluent material dispensing apparatus having a hopper for receiving and holding fluent material. As used herein, \"fluent material\" means material which flows or is capable of flowing and is made up of relatively small particles, such as powders, sand, gravel, rocks, pebbles, dirt, soil, limestone wastes, cement, grain, fertilizer or and other granular or powdery material.\nThis invention teaches several improvements on the parent application. In particular, the invention teaches a hopper having two opposing sides which converge at the lower end of the hopper where the multiple discharge chutes and the auger are located. The auger rotates to direct the filler material toward the four discharge chutes located at the bottom of the hopper. The rotation of the auger helps maintain an even flow of filler material through the discharge chutes. The invention also teaches the utilization of a means to drive the auger consisting of an engine or motor connected to a hydraulic pump and hydraulic motor system, and a means to control the three states of the auger: clockwise rotation, counter-clockwise rotation and no rotation.\nAlthough, other material handling devices are well known, such conventional equipment requires extensive material transmission apparatus and is not readily portable or self-contained. Furthermore, other material handling devices similar to the present invention do not teach a portable and efficient, power-driven mechanism for rotating an auger which moves fill material to the discharge openings so that blockages are eliminated and fill material flows constantly to the discharge openings. This invention overcomes the problems encountered by blockages often caused by moisture content within the fluent material. At the same time, this invention teaches a machine which is readily portable, having a self-contained auger power system.\nThe power system taught by this invention combines a conventional power source, such as a small-bore engine, with a hydraulic system, including a hydraulic pump and motor, connected by power transmission apparatus, to drive the auger. This relatively lightweight power system provides a material dispensing device that is easily portable on a trailer or in a truck bed, so that it may be taken on short notice to places where flooding occurs, or other form of material dispensing is needed.\nMaterial handling devices of a screw conveyor type are known, an example of a device of this type is illustrated in U.S. Pat. No. 3,093,271. However, that device has the shortcomings of requiring multiple augers necessary to discharge the filler material and also does not disclose a power-driven mechanism for rotating the auger. Similarly, French Patent No. 2641-262-A teaches a series of three augers adjacent to each discharge chute to reduce the risk of blockages. That device has the limitation that excessive numbers of augers are needed to discharge material from the hopper. Conversely, the present invention teaches a single auger which traverses the hopper and moves fill material to the multiple discharge chutes."}
{"text": "1. Field of the Invention\nThis invention relates generally to axisymmetric variable throat thrust vectoring nozzles and, more particularly, to the support apparatus of the vectoring ring that is used to pivot the nozzle flaps that vector the nozzle exhaust flow.\n2. Description of Related Art\nFor military aircraft applications, there exists a need to increase the maneuverability of the aircraft both for air to air combat missions and complicated ground attack missions. Aircraft designers seek to replace or augment the use of conventional aerodynamic surfaces such as flaps and ailerons with vectorable nozzles which turn or vector the exhaust flow and thrust of the gas turbine engine powering the aircraft. Hauer in U.S. Pat. No. 4,994,660, herein incorporated by reference, discloses an axisymmetric vectoring exhaust nozzle that provides a means for vectoring the thrust of an axisymmetric convergent/divergent nozzle by universally pivoting the divergent flaps of the nozzle in an asymmetric fashion or in other words pivoting the divergent flaps in radial and tangential directions with respect to the unvectored nozzle centerline. The flaps are pivoted by a vectoring ring which can be axially translated and gimballed or tilted about its horizontal and vertical axis (essentially have its attitude adjusted) through a limited range. Previously, two dimensional nozzles have used vectored thrust means that employ relatively flat flaps to direct the pitch or yaw direction of the engine's thrust.\nVectored thrust produces tangential and radial loads referred to as side loads that are transmitted from the flaps by various load paths back to the engine casing through the actuators. These tremendous loads require heavy actuators to absorb the loads and particularly the bending moments exerted on the actuator shafts by thrust vectoring.\nIt is therefore one object of the present invention to provide a means to minimize or eliminate the side loads transferred by the nozzle to the actuators. It is another object of the present invention to reduce or eliminate the bending moments that the actuators would be subject to due to the radial loads. Another object is to minimize the size and weight of the nozzle actuators and hydraulic system used to power the actuators.\nThese objects and other features and advantages will become more readily apparent in the following description when taken in conjunction with the appended drawings."}
{"text": "Portable ramps are used for loading railway maintenance equipment and other heavy machinery onto flat railcars that are to be transported to other service locations. These portable ramps may have multiple segments connected to one another with hinges to allow the segments to be folded for ease of transportation and storage of the ramp.\nWhen the ramps are deployed, heavy machinery can be driven or pulled up the ramp onto the flat railcar. Because of differences in the height of different parts of the rails on which the railcar rides, hinged ramps may sag in some areas, especially when the ramp is exposed to the heavy weight of machinery as it is loaded onto the railcar. This sagging causes the ramp to be uneven at the joints of the segments of the ramp.\nOne solution for supporting ramps at their hinges is described in U.S. Pat. No. 5,697,754 (the '754 patent). The '754 patent discloses a hinged ramp assembly for use with a flat railcar. To deploy and retract the ramp, a number of mechanical deployment and retraction facilitating elements are mounted to the ramp sections and operatively connected to a plurality of actuators. The '754 patent discloses legs connected to the ramp close to the hinged joints to provide additional support for sections of the ramp that are significantly higher than the surface of the rails.\nThe solution provided by the '754 patent is subject to a number of possible drawbacks. First, the stability of the assembly ramp in the '754 patent depends upon the ramp being deployed on a level surface. The '754 patent does not disclose any way to level the ramp when employed on uneven surfaces. For ramps deployed on rails, uneven rails can significantly affect the stability of the ramp, making uneven parts susceptible to sagging. Furthermore, the '754 patent only discloses using legs to support the hinges of those ramp sections closest to the railcar. It may be desirable to support the lower sections of the ramp, even where the ramp itself may be in contact with at least a part of one of the rails.\nThe presently disclosed system and method is directed to mitigating or overcoming one or more of the possible drawbacks set forth above and/or other problems in the art."}
{"text": "Consolidating workloads relates to allocating computing resources to multiple programs, for example assigning multiple programs to computing devices for execution, while attempting to consolidate the programs to be executed by only a part of the available devices. Such consolidation may enable, for example, the other devices to assume idle mode and reduce expenses by saving power or other resources.\nOne specific example of workload consolidation relates to executing virtual machines and attempting to reduce the number of used servers by consolidating the virtual machines.\nA virtual machine (VM) is a software implementation of a machine (i.e., a computer) that executes programs, similarly to a physical machine. There are two common types of virtual machines: 1. A system virtual machine which provides a complete system platform that supports execution of a complete operating system. A system virtual machine usually emulates an existing architecture, and may be used to provide a platform for executing programs where the real hardware is not available, for example, an obsolete or a very expensive platform. Another common use of system virtual machine relates to emulating multiple instances of identical machines, thus leading to more uniformity and efficient use of computing resources. 2. A process virtual machine, sometime referred to as language virtual machine, designed to support a single process and execute a single program. Process virtual machines are usually closely suited to one or more programming languages and built for providing program portability and flexibility.\nAn essential characteristic of a virtual machine of either type is that the software executed by the virtual machine is limited to the resources and abstractions provided by the virtual machine, and cannot exceed them even if the physical machine offers additional resources.\nOne purpose of consolidating virtual machines to physical machines is using fewer machines, thus saving power, maintenance and other resources. However, consolidation is to be performed such that each virtual machine existing on the target machine, as well as a virtual machine migrated to the physical machine will have sufficient resources such that it can execute as expected."}
{"text": "1. Field of the Invention\nThis invention relates to an improved digital signal processing method and apparatus for demultiplexing and multiplexing frequency division multiplexed channels and more particularly, to a lightweight, compact, low power system which synthesizes a polyphased filter, fast Fourier transform architecture suitable for use in applications with a large number of channels.\n2. Description of the Prior Art\nIn the prior art, digital polyphase filters (PPFs) are used in combination with fast Fourier transform (FFT) algorithms to efficiently demultiplex multiple, frequency division multiplexed (FDM) channels. Similarly, the polyphase filters are used in combination with inverse FFT algorithms to multiplex separate signals onto FDM channels. An N branch polyphase filter is realized by factoring (or partitioning) a desired prototype filter's transfer function (or impulse response) into N subfilters. When coupled with an N point FFT (or IFFT) as shown in FIGS. 1 and 2, an N branch PPF can be used to implement a bank of N bandpass, N:1 decimating (or 1:N interpolating) filters for demultiplexing (or multiplexing) up to N FDM channels. The MCD/M thus realized represents the most computationally efficient system known for both demultiplexing and multiplexing large numbers of FDM channels.\nThe minimum order necessary to synthesize the required polyphase prototype filter is dictated by the ratio of the guard frequency band (separating adjacent FDM channels) to channel center frequency spacing, filter passband gain uniformity, channel cross-talk attenuation requirements, and the number, N, of FDM channels to be either demultiplexed or multiplexed. Given a fixed guard band to channel spacing ratio and fixed filter response requirements, the minimum order of the polyphase prototype filter grows linearly with N.\nCellular, Personal Communication Systems (PCSs) and other emerging telecommunications applications continue to demand increasingly larger numbers of ever more closely packed FDM communications channels. Such systems (e.g., satellite based cellular systems) often place stringent constraints on the power, mass and physical dimensions of the digital MCD/M subsystem designed to generate the constituent FDM communications channels. These constraints necessitate the use of short sample word lengths (e.g., from 8 to 12 bits) and fixed point arithmetic in any practically realizable MCD/M system for these applications. Short sample word lengths and fixed point arithmetic are, however, inherently incompatible with accurate, high order prototype filter realizations. Indeed, in some applications, the indicated (higher order prototype) polyphase filter may not even be physically realizable within the finite word length and fixed point arithmetic dimensions necessary to meet the MCD/M subsystem's power, mass and size constraints."}
{"text": "1. Field of the Invention\nThe invention generally relates to medical devices and techniques, and more particularly to cardiovascular tissue closure devices and techniques.\n2. Description of the Related Art\nIn most cardiology and radiology procedures, a catheter is inserted into an artery, such as the femoral artery, through a vascular introducer. When the procedure is complete, the physician removes the catheter from the introducer and then removes the introducer from the arteriotomy in the vessel. The physician then must prevent or limit the amount of blood that leaks through the arteriotomy so that the patient can be discharged. Physicians currently use a number of methods to close the arteriotomy, such as localized compression, sutures, collagen plugs, and adhesives, gels, foams, clips, and similar materials.\nIn performing localized compression, the physician presses down against the vessel to allow the arteriotomy to naturally clot. This method, however, can take a significant amount of time, and requires the patient to remain immobilized and be kept in the hospital for observation. Clots at the puncture site may also be dislodged. Moreover, the amount of time necessary for the compression can significantly increase depending upon how much heparin, glycoprotein IIb/IIA antagonists, or other anti-clotting agents were used during the procedure. Sutures and collagen plugs can have procedure variability, can require time to close the vessel, and can necessitate a separate deployment device. Adhesives, gels, foams, and clips can have negative cost factors, can necessitate a complicated deployment process, and can have procedure variability."}
{"text": "In a voice packet transmission system, packets experience varying delays which depend on the transmission load at switching and multiplexing nodes and the number of such nodes in the packet transmission path. Heretofore, individual packets which experienced a delay greater than the so-called \"normal\" delay were considered late packets and were simply discarded. The discarding of such late packets causes distortion in a decoder which is attempting to reconstruct a replica of the transmitted speech signal. The distortion may be in the form of so-called clicks and pops in the reconstructed speech signal. This is extremely annoying and undesirable."}
{"text": "1. Field of the Invention\nThe present invention relates to a screen transmission device, screen display method and computer readable medium, and more particularly, to a window display position control technology in a multi-display system which displays a computer screen of a PC (Personal Computer) or the like by coordinating a plurality of displays.\n2. Related Art\nConventionally, as a system for displaying one computer screen of a PC or the like using a plurality of displays, there is a multi-monitor function supported as standard by Windows (registered trademark). When a display card for multi-monitor is used or a plurality of display cards are pushed into the PC simultaneously and the OS (Operating System) recognizes the plurality of displays, screens are displayed as one big desktop on the plurality of displays. The user can use a wider work space when using a plurality of displays than the case where the user is working only with one display. For example, when the user wants to see a plurality of documents simultaneously, document windows do not overlap with one another or windows are arranged so as not to overlap with one another, and therefore the windows can be displayed simultaneously without being reduced in size.\nA system whereby a plurality of displays cooperate with one another for displays is described in JP-A 2003-61140(KOKAI). According to JP-A 2003-61140(KOKAI), when a display instruction is received from a user, the system judges whether or not all the instructed display contents can be displayed on a full screen at a time. When all the instructed display contents cannot be displayed on the screen at a time, the system searches nearby devices and if it finds them, the system displays one screen using a plurality of devices coordinating with one another.\nFurthermore, the system described in JP-A 2004-318207(KOKAI) automatically collects information on screen creating elements and screen display elements included in a plurality of nearby devices and automatically forms a processing configuration which combines the screen creating elements and screen display elements in user specified order or in a preset priority order.\nWhile a multi-monitor function of Windows presupposes a relatively fixed way of combination between a PC and displays such as a desktop, using the methods described in the above described two patents facilitates the formation of a multi-display environment when a multi-display is necessary at a time of using a portable device or the like.\nAs described above, the systems described in JP-A 2003-61140(KOKAI) and JP-A 2004-318207(KOKAI) allow a multi-display environment to be easily used even when the apparatus configuration changes frequently such as when portable apparatuses are used. However, when a display is started up in a multi-display environment, which display is started up greatly affects the convenience of the user. For example, in a conventional Windows multi-monitor environment, a completely new window is started up at a primary monitor. When a display to be shown is not a primary monitor, it is necessary to move the window from the primary monitor to a desired display after the window is started up through mouse operation or the like. When the window is started up for the second time, it is started up at the position where it was closed last time, but in an apparatus configuration which changes fast from one moment to another, it would be rather unlikely that the position where the window was closed last time coincides with the desired position.\nThe system described in JP-A 2003-61140(KOKAI) can display one big screen using a plurality of devices in cooperation, but it does not refer to the position of a window when the window is started up. On the other hand, the system described in JP-A 2004-318207(KOKAI) simplifies or automates the determination of an apparatus configuration but it does not perform control on which apparatus a new window is displayed after the apparatus configuration is determined. In any case, there is a problem that a complicated operation to move the window to a desired display is required to display the window on a desired display when the window is started up."}
{"text": "The conventional approach for platform antenna isolation is to place the antenna component at essentially the farthest distance from the other components of the platform in order to obtain minimum undesired antenna coupling. This is limited by a compact platform with limited space. Currently there is a need for antenna isolation in multi-radio platforms, particularly those that employ multiple input-multiple outputs (MIMO)."}
{"text": "Compression bandages are frequently used in medical and sports applications requiring a strong and reliable, yet comfortable and easily applied, means of securing a limb or other body segment for prolonged periods of time. For example, strains and sprains can cause inflammation and the accompanying accumulation of fluid around a sprained joint. Wrapping the affected joint securely with an elastic bandage can prevent excess fluid from accumulating and causing additional tissue damage.\nIn addition, chronic venous disease, including valve insufficiency and venous wall damage, and leg ulcers of various origins, including venous stasis ulcers, arterial (ischemic) ulcers and neurotrophic ulcers, are common medical problems. Leg ulcers are wounds or open sores that do not heal, or otherwise recur repeatedly, and cause persistent swelling as well as burning, itching, irritation and discoloration of the skin. Treatment of leg ulcers generally includes compression in addition to topical protection of the wound and antimicrobial treatment of the affected area.\nTherapeutic compressive pressures are based on the principle that extra-vascular pressure should equal excess venous pressure in order to restore normal venous flow from the extremities back to the torso, e.g., from the foot to the thigh. Stress on the venous system is already the greatest in the leg in standing position because the veins must move blood against the force of gravity. Ideal leg compression mirrors the leg's natural degressive pressure gradient, i.e., higher pressure at the ankle and lower pressure at the thigh, which results from the gradual increase of leg diameter from the ankle to the thigh. Thus, therapeutic compressive pressure goals vary from patient to patient and at different locations in each patient's leg based on venous pressure measurements.\nTraditionally, compression therapy for leg ulcers has involved wrapping a patient's leg with a layer of cotton batting followed by a compression bandage or sleeve. Better results are now achieved with the application of a 3 to 4 layer bandage system, by which different types of bandages are combined to achieve and sustain a suitable therapeutic sub-bandage compressive pressure. A common example is the four-layer system described in “The Function of Multiple Layer Compression Bandaging in the Management of Venous Ulcers,” D D I Wright et al., SWM, 10, 109-10 (1988). The four layers include (1) a skin-facing layer of cotton wool, (2) a lightweight bandage, (3) a light compression bandage, and (4) a flexible cohesive bandage.\nDespite their potential to deliver therapeutic compressive pressure, the process of applying such 3 to 4 layer bandage systems is time intensive, which reduces patient compliance. Also, the bandages in these systems often require pleats or tucks in order to conform to the patient's leg and may slip or wrinkle following application, all of which can irritate the underlying skin and make the results of compression highly variable. Patient comfort and compliance is further reduced by the high profile or thickness and lack of thermal and/or moisture regulation of 3 to 4 layer systems, which can be difficult to wear under existing clothing or footwear and cause the uncomfortable buildup of body heat and moisture under the bandages.\nAlternatively, other compression bandage systems have been proposed in attempts to counteract the complexities and time demands involved with the application of 3 to 4 layer bandage systems. For example, U.S. Pat. No. 7,854,716 discloses a two-part compression bandage system. The first part is an inner multi-layer elastic bandage comprising (1) an elastic substrate layer with a self-adhering outer face and an inner face affixed to (2) a foam layer having a skin-facing exposed face not affixed to the substrate layer. The second part of the system is a separate outer elastic layer with self-adhesion and compression properties. In application, the inner multi-layer bandage is wrapped around the patient's leg with the elastic substrate outer face configured and adapted to adhere to the separate outer elastic layer to prevent slippage or migration. Thus, the only purpose of the expensive elastic substrate layer of the inner bandage is to facilitate adhesion with the separate outer elastic layer. In addition to cost concerns, the two-part system falls short of optimizing patient comfort and, as a result, patient compliance with regard to its thickness, conformability, heat retention, and moisture regulation.\nMoisture regulation in a compression bandage system is particularly important due to the seeping of serous fluid and other wound transudate and/or exudate from leg ulcers or other chronic wounds. Traditionally, cotton batting has been used under compression bandages for fluid absorption; however, the cotton often sticks to the affected area upon removal. Also of concern is the ability of any sub-bandage absorbent padding to remain compression neutral and minimize friction against the skin."}
{"text": "The present invention relates to an automatic focus control device having a plurality of focus detecting regions.\nConventionally, an automatic focus detecting device having a single focus detecting region, in which a defocus amount is corrected in accordance with assembly errors of an optical system, has been proposed in, for example, Japanese Patent Laid-Open Publication No. 126517/1984. However, an automatic focus detecting device having a plurality of focus detecting regions, in which a defocus amount is corrected for each of the focus detecting regions, is not known.\nAs shown in FIG. 7, in the case where automatic focus detection is performed in three focus detecting regions, i.e. a horizontally elongated focus detecting area A disposed at a center of a photographing field S and a pair of vertically elongated focus detecting regions B and C disposed at opposite sides of the photographing field S, CCD light-receiving element arrays e1, e2 and e3 corresponding to the focus detection regions A, B and C, respectively are provided on one chip as shown in FIG. 4. If only one CCD light-receiving element array, for example, only the CCD light-receiving element array corresponding to the central focus detecting region A is provided, positional adjustment of the CCD chip in the direction of its optical axis can be performed by performing positional adjustment of only the central portion of the CCD chip in the direction of the optical axis. However, in the case where a plurality of, for example, three CCD light-receiving element arrays are provided as described above, parallelism of the CCD chip plays an important role. Namely, when parallelism of the CCD chip is poor, deviation in the direction of the optical axis occurs at opposite sides of the CCD chip. When parallelism of the CCD chip is poor, deviation in the direction of the optical axis occurs at opposite sides of the CCD chip. Even in the case of only one CCD light-receiving element array, the deviation also takes place at opposite sides of the CCD chip but assumes a small value. Meanwhile, also in the case where a plurality of the CCD light-receiving element arrays corresponding to the focus detecting regions, respectively are not provided on a single chip, deviations of the assembled CCD light-receiving element arrays in the direction of the optical axis are different from each other, so that it becomes necessary to correct the respective focus detecting regions in the direction of the optical axis.\nAs shown in FIG. 6a, a focus detecting device has been widely used in which an image formed by a light-receiving lens a is re-formed by a pair of reimaging lenses d1 into first and second images on a light.TM.receiving element array el arranged in a straight line and a distance between the first and second images is detected by the light-receiving element array el such that a focusing state is detected. In this known focus detecting device, when the, distance between the first and second images assumes a predetermined value, a decision of an in-focus state in which the image is formed on a predetermined focal plane is made. Meanwhile, when the distance between the first and second images is smaller and larger than the predetermined value, decisions of front and rear focus states in which the image is formed forwardly and rearwardly of the predetermined focal plane are made, respectively such that deviation amounts of the distance between the images from that of the in-focus position are outputted as defocus amounts. At the time of focus detection, the distance j between the first and second images is obtained and the defocus amount .DELTA..epsilon. is calculated by multiplying the distance j by a predetermined coefficient s. The above described Japanese Patent Laid-Open Publication proposes correction of the coefficient s since the coefficient s for the front focus state is different from that for the rear focus state.\nIn the CCD chip of FIG. 4, it is considered that if a length of the CCD light-receiving element arrays e1 and e3 corresponding to the focus detecting regions B and C disposed at opposite sides of the photographing field, respectively is made smaller than that of the CCD light-receiving element array e2 corresponding to the focus detecting region A disposed at the center of the photographing field, a focus detecting module can be made smaller in size by a difference between the length of the CCD light-receiving element arrays e1 and e3 and that of the CCD light-receiving element array e2 and a time period required for effecting data dump from the CCD light-receiving element arrays el and e3 can be reduced. To this end, as shown in FIG. 4, a pair of reimaging lenses d1 and a pair of reimaging lenses d3, which are, respectively, disposed at opposite sides of the CCD chip, have distances Dd1 and Dd3 smaller than a distance Dd2 between a pair of reimaging lenses d2 disposed at a center of the CCD chip, so that the numbers of the elements required for the CCD light-receiving element arrays e1 and e3 are determined accordingly. However, in this case, such a problem arises due to the difference between the distances Dd1 and Dd3 and the distance Dd2 that a coefficient for converting a distance between the first and second images into a defocus amount on the CCD light-receiving element array e2 is different from a coefficient for converting a distance between the first and second images into a defocus amount on the CCD light-receiving element array e1 and a coefficient for converting a distance between the first and second images into a defocus amount on the light-receiving element array e3.\nIn FIG. 6a, the distance De1 between the first and second images on the light-receiving element array e1 changes according to changes of a distance Dd1 of the reimaging lenses d1. Therefore, in the case where the reimaging lenses d1 are made of, for example, plastic, the distance Dd1 of the reimaging lenses d1 varies in response to temperature change, so that the distance De1 between the first and second images on the light-receiving element array e1 changes, thereby resulting in also change of the defocus amount. Thus, Japanese Patent Laid-Open Publication No. 235110/1985 discloses that temperature of the focus detecting device is detected such that the defocus amount is corrected in accordance with the detected temperature of the focus detecting device.\nMeanwhile, if temperature rise occurs in the case where as shown in FIG. 4, a pair of the reimaging lenses d1 and a pair of the reimaging lenses d3, which are, respectively, disposed at opposite sides of the CCD chip, have the distances Dd1 and Dd3 smaller than a distance Dd2 of a pair of the reimaging lenses d2 disposed at a center of the CCD chip, so that the numbers of the elements required for the CCD light-receiving element arrays e1 and e3 are determined accordingly, amounts of change of the distances between the two images on the CCD light-receiving element arrays e1, e2 and e3 upon temperature rise are different from each other due to the differences among the distances Dd1, Dd2 and Dd3 even if the reimaging lenses d1, d2 and d3 are made of an identical material. Meanwhile, in the case where the reimaging lenses are not made of an identical material, the reimaging lenses have different coefficients of thermal expansion, so that amounts of change of the distances between the images on the CCD light-receiving element arrays e1, e2 and e3 upon temperature rise are different from each other for the respective focus detecting regions."}
{"text": "U.S. Pat. Nos. 4,017,615, 3,888,995 and 3,592,930 pertain to ointment-like vehicle for corticosteroid drugs and the like.\nReferences which disclose film formers as ingredients for topically administered dry formulations include U.S. Pat. Nos. 2,973,300, 3,214,338 and 4,210,633.\nBandage-like devices for delivering drugs topically or transdermally are described in U.S. Pat. Nos. 3,598,123, 4,191,743, 4,605,548, 4,655,768, 4,409,206, 4,286,592, 4,230,105, 3,948,254, 3,742,951, 3,734,097 and 3,731,683. Such devices may include separate adhesion and drug reservoir layers.\nU.S. Pat. No. 3,536,809 discloses a mixture of a drug such as progesterone dispersed in a polyalkyleneglycol impregnated into a fabric strip which may be retained in the mouth to administer the drug through the buccal mucosa.\nConsidered more pertinent to the invention hereof are polymeric dispersion matrix materials which have skin adhesive properties and have a drug dispersed directed into the matrix.\nIn U.S. Pat. No. 3,632,740 it is suggested to incorporate a drug such as a corticosteroid into a pressure sensitive rubbery adhesive layer on a flexible backing.\nIn U.S. Pat. No. 4,292,301 there is disclosed a polymeric diffusion matrix said to permit sustained release of ephedrine, the matrix comprising a polar plasticizer such as glycerol or polyethylene glycol (MW 1000), polyvinyl alcohol, polyvinyl pyrrolidone and ephedrine.\nU.S. Pat. No. 4,470,962 discloses a polymeric diffusion matrix said to be capable of sustained release of a drug comprising glycerol, polyvinylalcohol, a water soluble polymer with hydration sites, a drug dispersed therein and water.\nU.S. Pat. Nos. 4,307,717 and 4,675,009 describe bandage materials which include a backing element and a substrate comprising a matrix material comprising a solid phase formed of a polysaccharide or certain synthetic polymers and a liquid phase consisting of a hydric alcohol, carbohydrates and/or proteins in an aqueous solution. The matrix material also contains a medicament suspended or dissolved therein.\nU.S. Pat. No. 4,593,053 discloses polyvinylalcohol/polyvinyl pyrrolidone based gels with non-syneresing adhesive characteristics. In one embodiment the gels contain an ionic drug which can be iontophoretically delivered. Tacifiers such as poly-2-acrylamido-2-methylpropane sulfonic acid may be added to the formulation at levels of 2%-20%.\nHydrogel materials based on polymerized 2-acrylamido-2-methylpropane sulfonate salts are known as electrode materials from U.S. Pat. Nos. 4,391,278 and 4,768,523. In U.S. Pat. No. 4,391,278 a tape electrode is disclosed which includes a gel comprising polymerized 2-acrylamido-2-methane sulfonic acid or a salt thereof, water and/or an alcohol to give electrically conductive, flexible skin adhering properties. All of the example formulations employing an alcohol utilize glycerol although propylene glycol and sorbitol are also mentioned as useful."}
{"text": "This invention pertains to a cellulosic fiberboard panel, which is a building construction panel reinforced in a novel manner. This invention contemplates that reinforcing bands, which may be cementitious, are applied on an expansive surface of a cellulosic fiberboard panel.\nCommonly, cellulosic fiberboard panels are made by severing each panel, via water jets or otherwise, from a continuous web of cellulosic fiberboard, as made on a fiberboard-making machine, such as a Fourdrinier machine, from cellulosic fibers, which may be wood fibers, bagasse fibers, or other cellulosic fibers. Cellulosic fiberboard panels are used widely in building construction, as wall sheathing, as insulating panels, as roofing underlayments, as flooring underlayments, and otherwise.\nCharacteristically, unless made, coated, or treated with cementitious materials, asphaltic materials, or other non-permeable materials, cellulosic fiberboard panels are permeable to water vapor. Because of their permeability to water vapor, cellulosic fiberboard panels are desirable for installations where mold or mildew might thrive unless water vapor that would be otherwise trapped could permeate such panels, such as installations of wall sheathing over spaced studs, at outer walls of buildings. However, as compared to plywood panels and to other building panels, cellulosic fibreboard panels are not as strong, unless reinforced.\nAs exemplified in numerous prior patents including U.S. patent Re. 32,329, which is a reissue of U.S. Pat. No. 4,339,405, and including U.S. Pat. Nos. 4,402,751, No. 4,406,703, No. 4,981,518, No. 5,744,078, No. 5,846,894, and No. 6,001,169, it has been known heretofore to blend cementitious materials with cellulosic fibers, whereby to make reinforced, water vapor-impermeable, cellulosic fiberboard panels.\nIn a building construction panel made of cellulosic fiberboard and having two expansive surfaces, two longitudinal edges, and two transverse edges, each expansive surface having four corners and each expansive surface having an overall area, this invention provides an improvement wherein the panel has bands of a reinforcing material, for which a cementitious material is preferred. The bands cover a minor portion of the overall area of a given one of the expansive surfaces. The panel is devoid of the reinforcing material except at the bands.\nThe bands may include two diagonal bands, each of which extends between diagonally opposite ones of said corners of the given one of the expansive surfaces. The bands may include two longitudinal bands, each of which extends along a respective one of the longitudinal edges of the given one of the expansive surfaces. The bands may include two transverse bands, each of which extends along a respective one of the transverse edges of the given one of the expansive surfaces. Preferably, the bands include such diagonal, longitudinal, and transverse bands.\nAs a method to reinforce a panel made of cellulosic fiberboard and having two expansive surfaces, two longitudinal edges, and two transverse edges, each expansive surface having four corners and each expansive surface having an overall area, this invention provides a step of wherein applying bands of a reinforcing material, for which a cementitious material is preferred, so as to cover a minor portion of the overall area of a given one of the expansive surfaces, wherein the panel is devoid of the reinforcing material except at the bands.\nTwo diagonal bands of the reinforcing material may be thus applied so that each of the diagonal bands extends between diagonally opposite ones of said corners of the given one of the expansive surfaces. Two longitudinal bands of the reinforcing material may be thus applied so that each of the longitudinal bands extends along a respective one of the longitudinal edges of the given one of the expansive surfaces. Two transverse bands of the reinforcing material may be thus applied so that each of the transverse bands extends along a respective one of the transverse edges of the given one of the expansive surfaces. Preferably, such diagonal, longitudinal, and transverse bands are applied so as to extend as noted."}
{"text": "In recent years, there has been significant technological progress in the field of displays. In particular, the needs of the market have stimulated tremendous progress in the technology directed to increasing resolution and thinning displays.\nIn the next phase of this field, focus is placed on commercialization of flexible displays having a curved display area, and a variety of proposals have been made to increase the flexibility of displays (for example, see Patent Document 1). In addition, light-emitting devices using a flexible substrate can be lightweight compared to those using a substrate of usual thickness which is formed of glass or the like.\nAs methods for manufacturing a light-emitting device using a flexible substrate, the following proposals have been made: a method in which a light-emitting element and other elements are directly provided over a flexible substrate; a method in which a light-emitting element is formed over a glass substrate of usual thickness and then the substrate is subjected to polishing treatment or the like, so that the substrate is thinned to have flexibility or the substrate is removed and the light-emitting element is attached to a flexible substrate; a method in which a light-emitting element and other elements are formed over a glass substrate of usual thickness and then a layer having the light-emitting element and the other elements is separated from the substrate and transferred to a flexible substrate; and the like.\nNevertheless, over a substrate having sufficient flexibility, a light-emitting element and other elements are not easy to form with high precision. Therefore, over such a substrate, it is difficult to form light-emitting elements for their respective colors or to provide a color filter, much less to form a semiconductor element. Further, a light-emitting device using a thinned glass substrate can be said to lack sufficient flexibility. Even if a light-emitting device is formed by removal of a glass substrate, such a light-emitting device has a problem of low productivity. In contrast, the method for manufacturing a light-emitting device in which a separation step and a transfer step are utilized is relatively simple and easy and facilitates fabrication of a semiconductor element or formation of light-emitting elements for their respective colors. This method can also ensure sufficient flexibility, which shows great promise."}
{"text": "Rotating wing aircraft rely on a rotating wing to provide lift. In contrast, fixed wing aircraft rely on air flow over a fixed wing to provide lift. Fixed wing aircraft must therefore achieve a minimum ground velocity on takeoff before the lift on the wing is sufficient to overcome the weight of the plane. Fixed wing aircraft therefore generally require a long runway along which to accelerate to achieve this minimum velocity and takeoff.\nIn contrast, rotating wing aircraft can take off and land vertically or along short runways inasmuch as powered rotation of the rotating wing provides the needed lift. This makes rotating wing aircraft particularly useful for landing in urban locations or undeveloped areas without a proper runway.\nThe most common rotating wing aircraft in use today are helicopters. A helicopter typically includes a fuselage, housing an engine and passenger compartment, and a rotor, driven by the engine, to provide lift. Forced rotation of the rotor causes a reactive torque on the fuselage. Accordingly, conventional helicopters require either two counter rotating rotors or a tail rotor in order to counteract this reactive torque.\nAnother type of rotating wing aircraft is the autogyro. An autogyro aircraft derives lift from an unpowered, freely rotating rotor or plurality of rotary blades. The energy to rotate the rotor results from a windmill-like effect of air passing through the underside of the rotor. The forward movement of the aircraft comes in response to a thrusting engine such as a motor driven propeller mounted fore or aft.\nDuring the developing years of aviation aircraft, autogyro aircraft were proposed to avoid the problem of aircraft stalling in flight and to reduce the need for runways. The relative airspeed of the rotating wing is independent of the forward airspeed of the autogyro, allowing slow ground speed for takeoff and landing, and safety in slow-speed flight. Engines may be tractor-mounted on the front of an autogyro or pusher-mounted on the rear of the autogyro.\nAirflow passing the rotary wing, alternately called rotor blades, which are tilted upward toward the front of the autogyro, act somewhat like a windmill to provide the driving force to rotate the wing, i.e. autorotation of the rotor. The Bernoulli effect of the airflow moving over the rotor surface creates lift.\nVarious autogyro devices in the past have provided some means to begin rotation of the rotor prior to takeoff, thus further minimizing the takeoff distance down a runway. One type of autogyro is the “gyrodyne,” which includes a gyrodyne built by Fairey aviation in 1962 and the XV-1 convertiplane first flight tested in 1954. The gyrodyne includes a thrust source providing thrust in a flight direction and a large rotor for providing autorotating lift at cruising speeds. To provide initial rotation of the rotor, jet engines were secured to the tip of each blade of the rotor and powered during takeoff, landing, and hovering.\nAlthough rotating wing aircraft provide the significant advantage of vertical takeoff and landing (VTOL), they are much more limited in their maximum flight speed than are fixed wing aircraft. The primary reason that prior rotating wing aircraft are unable to achieve high flight speed is a phenomenon known as “retreating blade stall.” As the fuselage of the rotating wing aircraft moves in a flight direction, rotation of the rotor causes each blade thereof to be either “advancing” or “retreating.”\nThat is, in a fixed-wing aircraft, all wings move forward in fixed relation, with the fuselage. In a rotary-wing aircraft, the fuselage moves forward with respect to the air. However, rotor blades on both sides move with respect to the fuselage. Thus, the velocity of any point on any blade is the velocity of that point, with respect to the fuselage, plus the velocity of the fuselage. A blade is advancing if it is moving in the same direction as the flight direction. A blade is retreating if it is moving opposite the flight direction.\nThe rotor blades are airfoils that provide lift that depends on the speed of air flow thereover. The advancing blade therefore experiences much greater lift than the retreating blade. One technical solutions to this problem is that the blades of the rotors are allowed to “flap.” That is, the advancing blade is allowed to fly or flap upward in response to the increased air speed thereover such that its blade angle of attack is reduced. This reduces the lift exerted on the blade. The retreating blade experiences less air speed and tends to fly or flap downward such that its blade angle of attack is increased, which increases the lift exerted on the blade.\nFlap enables rotating wing aircraft to travel in a direction perpendicular to the axis of rotation of the rotor. However, lift equalization due to flapping is limited by a phenomenon known as “retreating blade stall.” As noted above, flapping of the rotor blades increases the angle of attack of the retreating blade. However, at certain higher speeds, the increase in the blade angle of attack required to equalize lift on the advancing and retreating blades results in loss of lift (stalling) of the retreating blade.\nA second limit on the speed of rotating wing aircraft is the drag at the tips of the rotor. The tip of the advancing blade is moving at a speed equal to the speed of the aircraft and relative to the air, plus the speed of the tip of the blade with respect to the aircraft. That is equal to the sum of the flight speed of the rotating wing aircraft plus the product of the length of the blade and the angular velocity of the rotor. In helicopters, the rotor is forced to rotate in order to provide both upward lift and thrust in the direction of flight. Increasing the speed of a helicopter therefore increases the air speed at the rotor or blade tip, both because of the increased flight speed and the increased angular velocity of the rotors required to provide supporting thrust.\nThe air speed over the tip of the advancing blade can therefore exceed the speed of sound even though the flight speed is actually much less. As the air speed over the tip approaches the speed of sound, the drag on the blade becomes greater than the engine can overcome. In autogyro aircraft, the tips of the advancing blades are also subject to this increased drag, even for flight speeds much lower than the speed of sound. The tip speed for an autogyro is typically smaller than that of a helicopter, for a given airspeed, since the rotor is not driven. Nevertheless, the same drag increase occurs eventually.\nA third limit on the speed of rotating wing aircraft is reverse air flow over the retreating blade. As noted above, the retreating blade is traveling opposite the flight direction with respect to the fuselage. At certain high speeds, portions of the retreating blade are moving rearward, with respect to the fuselage, slower than the flight speed of the fuselage. Accordingly, the direction of air flow over these portions of the retreating blade is reversed from that typically designed to generate positive lift. Air flow may instead generate a negative lift, or downward force, on the retreating blade. For example, if the blade angle of attack is upward with respect to wind velocity, but wind is moving over the wing in a reverse direction, the blade may experience negative lift.\nThe ratio of the maximum air speed of a rotating wing aircraft to the maximum air speed of the tips of the rotor blades is known as the “advance ratio. The maximum advance ratio of rotary wing aircraft available today is less than 0.5, which generally limits the top flight speed of rotary wing aircraft to less than 200 miles per hour (mph). For most helicopters, that maximum achievable advance ratio is between about 0.3 and 0.4.\nIn view of the foregoing, it would be an advancement in the art to provide a rotating wing aircraft capable of vertical takeoff and landing and flight speeds in excess of 200 mph."}
{"text": "1. Technical field\nThe present invention relates to a remote control system for controlling home appliances from outside the home and a gateway apparatus used for the remote control system, and in particular, to the system and apparatus capable of controlling the home appliances from outside the home with easier operations.\n2. Related art\nIn recent years, in order to build a system in which home appliances are connected to a network built inside the home (home network) so as to be controlled those appliances, standardization for protocols of the home network that use IEEE 1394, BLUETOOTH, and/or power lines is now in progress. The home appliances have already been partly mutual-connected for linked usage, and home appliances can be connected to personal computers or to other home appliances (such that a floor type of VTR is connected to a video camera).\nA Japanese Patent publication No. 3016350 discloses a system for a home appliance control interface, in which there is provided a personal computer for operating home appliances and an image indicative of rooms in which the home appliances are placed.\nAnother Japanese Patent publication No. 3014532 discloses a home automation system. In this system, the layout view of rooms of a house is displayed on a personal computer used for controlling home appliances, such as a bath and air conditioner. For example, when a user performs operations on the layout view of the air conditioner, contents of operations for a wind direction, wind amount, setting of temperature, necessity for dehumiditification and/or others are displayed on the view. Based on a selected content to be operated, the system is able to actuate the air conditioner. In this home automation system, the personal computer is connected to a second personal computer through a telephone line, so that a user is able to give data to the second personal computer to control the home appliances from a distance.\nWith regard to a control interface for home appliances, remote control devices have long been used. Japanese Patent publication 6-32509 discloses a remote control apparatus, in which replacement of ROMs allows a plurality of appliances to be controlled. This remote control apparatus has a display unit for displaying figures illustrating operation panels of the appliances to be controlled. Pressing a key on an operation panel that has been displayed enables control of a desired appliance.\nFurther, the Internet has widely become popular in recent years, and use of a home computer and portable terminals or mobile terminals makes it possible to have access to the Internet from any place. To control a personal computer with functions for home appliances (e.g., a personal computer with the functions for a TV tuner and picture recording) via the Internet has already been conducted. Practically, a portable terminal is used outside to send to a mail box on the Internet an electronic mail in which pieces of information indicating instructions to control a home appliance are written. On the other hand, a home personal computer is actuated to regularly access the electronic mail box. If the home personal computer finds an electronic mail containing the information in relation to instructions of control, the function of the personal computer is activated.\nIn this configuration, the portable terminal extracts a function, operation data for reserving picture recording, and a start time and an end time of recording of a broadcast channel. And the portable terminal uses these pieces of information to produce the information about the control instructions that can be interpreted by a program installed in the personal computer. The produced information is then sent out with the electronic mail.\nFor a further example of accessing a computer system via a communication network, a system has been known in which terminals control a host computer based on programs on API (application programming interface) specifications that include communication protocol processing.\nHowever, the conventional systems have given no consideration to a direct connection of a home network to the Internet, because there are differences in protocols between the home network connected to home appliances and an outside computer network (e.g., Internet) connected to the home network. In addition, the home appliances are not equipped with the API capable of receiving control commands sent from outside the home via the Internet and operating the home appliances. Therefore, there is a problem that the conventional systems are unable to remote-control home appliances from outside the home via the Internet, though the Internet itself has now considerably become popular.\nMoreover, in the case of conventional remote-control techniques, there is a possibility that conflicts occur between direct operations and remote operations toward a certain home appliance.\nStill further, there is a fear that an immoral user might perform a remote control to give home appliances a breakdown or information about the leakage of privacy."}
{"text": "Peptidyl-2-amino-1-hydroxyalkanesulfonic acid inhibitors of cysteine proteases, methods for making these, compounds, and methods for using the same are disclosed.\nNumerous cysteine proteases have been identified in human tissues. A xe2x80x9cproteasexe2x80x9d is an enzyme which degrades proteins into smaller components (peptides). The term xe2x80x9ccysteine proteasexe2x80x9d refers to proteases which are distinguished by presence therein of a cysteine residue which plays a critical role in the catalytic process. Mammalian systems, including humans, normally degrade and process proteins via a variety of enzymes including cysteine proteases. However, when present at elevated levels or when abnormally activated, cysteine proteases are involved in pathophysiological processes.\nFor example, calcium-activated neutral proteases (xe2x80x9ccalpainsxe2x80x9d) comprise a family of intracellular cysteine proteases which are ubiquitously expressed in mammalian tissues. Two major calpains have been identified; calpain I and calpain II. While calpain II is the predominant form in many tissues, calpain I is thought to be the predominant form in pathological conditions of nerve tissues. The calpain family of cysteine proteases has been implicated in many diseases and disorders, including neurodegeneration, stroke, Alzheimer\"\"s disease, amyotrophy, motor neuron damage, spinal cord trauma, traumatic brain injury, acute central nervous system injury, muscular dystrophy, bone resorption, platelet aggregation, cataracts and inflammation. Calpain I has been implicated in excitatory amino-acid induced neurotoxicity disorders including ischemia, spinal cord trauma, traumatic brain injury, hypoglycemia and epilepsy.\nThe lysosomal cysteine protease cathepsin B has been implicated in arthritis, inflammation, myocardial infarction, tumor metastasis, and muscular dystrophy. Other lysosomal cysteine proteases include cathepsins C, H, L and S. Interleukin-1xcex2 converting enzyme (ICE) is a cysteine protease which catalyzes the formation of interleukin-1xcex2. Interleukin-1xcex2 is an immunoregulatory protein implicated in inflammation, diabetes, septic shock, rheumatoid arthritis, and Alzheimer\"\"s disease. ICE has also been linked to apoptotic cell death of neurons which is implicated in a variety of neurodegenerative disorders including Parkinson\"\"s disease, ischemia, and amyotrophic lateral sclerosis (ALS).\nCysteine proteases are also produced by various pathogens. The cysteine protease clostripain is produced by Clostridium histolyticum. Other proteases are produced by Trpanosoma cruzi, malaria parasites Plasmodium falciparum and P.vinckei and Streptocococcus. Viral cysteine proteases such as HAV C3 are essential for processing of picornavirus structural proteins and enzymes.\nGiven the link between cysteine proteases and various debilitating disorders, compounds which inhibit these proteases would be useful and would provide an advance in both research and clinical settings.\nThe present invention is directed to novel cysteine protease inhibitors which we refer to as peptidyl-2-amino-1-hydroxyalkanesulfonic acids. Exemplary compounds are represented by Formula I: \nConstituent members and preferred embodiments are defined infra.\nThe compounds of the invention are useful for the inhibition of cysteine proteases. Beneficially, these compounds find utility in a variety of settings. For example, in a research arena, the subject compounds can be used as standards for screening in the discovery of agents for treating disorders associated with abnormal and/or aberrant activity of cysteine proteases. In a therapeutic arena, the compounds can be used to alleviate, mediate, reduce, and/or prevent disorders which are associated with abnormal and/or aberrant activity of cysteine proteases. One particular advantage of the subject compounds is that they have unexpectedly excellent solubility in aqueous media, approximately between 1 and 300 mg/mL. Thus, the compounds can be easily administered parenterally, for example intravenously, intrathecally or supradurally by bolus or infusion in the clinical or laboratory setting in media such as aqueous saline buffer.\nAlso disclosed are methodologies for making the peptidyl-2-amino-1-hydroxyalkanesulfonic acids. These and other features of the invention are set forth in more detail below.\nNovel cysteine protease inhibitors have been discovered which are represented by Formula I: \nwherein:\n* denotes the xcex1-carbon of an xcex1-amino acid residue having the L configuration;\no denotes a carbon having either stereochemical configuration, or a mixture thereof;\nA is selected from the group consisting of lower alkyl, aryl having from 6 to about 14 carbons, heterocyclyl having from about 5 to about 14 ring atoms, heterocycloalkyl having from about 5 to about 14 ring atoms, aralkyl having from about 7 to about 15 carbons, and heteroarylalkyl, said alkyl, aryl, heterocyclyl, heterocycloalkyl, aralkyl and heteroarylalkyl groups being optionally substituted with J;\nB is selected from the group consisting of C(xe2x95x90O), OC(xe2x95x90O) S(xe2x95x90O) S (xe2x95x90O)2, and NR4C(xe2x95x90O), where R4 is H or lower alkyl;\neach Aaa is independently an amino acid which optionally contains one or more blocking groups;\nn is 0, 1, 2, or 3;\nG is selected from the group consisting of H, C(xe2x95x90O)NR5R6, C(xe2x95x90O)OR5, CF3, CF2R5, P(xe2x95x90O)(R5)(OR6) and P(xe2x95x90O)(OR5)(OR6);\nJ is selected from the group consisting of halogen, lower alkyl, cycloalkyl, aryl, heteroaryl, heterocycloalkyl, aryl substituted with aralkyloxy, C(xe2x95x90O)OR7, OC(xe2x95x90O)R7, NR8C(xe2x95x90O)OR7, OR7, CN, NO2, NR7R8, Nxe2x95x90C(R7)R8, SR7, S(xe2x95x90O)R7, S(xe2x95x90O)2R7, and C(xe2x95x90NR7)NHR8;\nR5 and R6 are independently selected from the group consisting of H, lower alkyl, aralkyl, heterocyclic, and heterocycloalkyl, said lower alkyl, aralkyl, heterocyclic, and heterocycloalkyl groups being optionally substituted with one or more hydroxy, alkoxy, aryloxy, carboxy, alkoxycarbonyl, aryloxycarbonyl, amino, monoalkylamino, dialkylamino, monoarylamino, diarylamino, or halogen groups;\nR1, R2, and R3, independently, are selected from the group consisting of H, lower alkyl, aryl, and heterocyclyl, said lower alkyl, aryl and heterocyclyl groups being optionally substituted with one or more J groups;\nor R2 and R3, may be taken together with the carbon and nitrogen atoms to which they are attached to form a 4-8 membered ring which is optionally substituted with one or more J groups;\nR7 and R8, independently, are selected from the group consisting of H, lower alkyl, aralkyl, heterocyclic, and heterocycloalkyl, said lower alkyl, aralkyl, heterocyclic, and heterocycloalkyl groups being optionally substituted with one or more hydroxy, alkoxy, aryloxy, carboxy, alkoxycarbonyl, aryloxycarbonyl, amino, monoalkylamino, dialkylamino, monoarylamino, diarylamino, or halogen groups;\nM is a pharmaceutically acceptable cation selected from the group consisting of sodium, lithium, potassium, calcium, magnesium, zinc, aluminum, ammonium, mono-, di-, tri-, or tetraalkylammonium, morpholinium, piperidinium, and megluminium; and\nwith the proviso that R2 and R3 taken together is other than xe2x80x94CH2xe2x80x94CH2xe2x80x94CH2xe2x80x94.\nIn some preferred embodiments of the compounds of Formula I R1 is lower alkyl or lower alkyl substituted with J.\nIn some preferred embodiments, J is cycloalkyl, aryl, or aryl substituted with aralkyloxy. In more preferred embodiments, R1 is ethyl, isopropyl, benzyl, isobutyl, cyclohexylmethyl, or 4-benzyloxybenzyl.\nIn other preferred embodiments, R2 is alkyl or cycloalkyl. In more preferred embodiments, R2 is isobutyl, isopropyl, or cyclopentyl.\nIn some preferred embodiments, R3 is H. In other preferred embodiments, R2 and R3 are taken together with the carbon and nitrogen atoms to which they are attached to form a 4-8 membered ring which is optionally substituted with one or more J groups.\nIn some preferred embodiments, G is H or C(xe2x95x90O)NHEt. More preferably, G is H.\nIn further preferred embodiments, A is alkyl, aryl, aralkyl, or tetrahydroisoquinolin-2-yl. In especially preferred embodiments, A is methyl, tolyl, naphthyl, benzyl, or tetrahydroisoquinolin-2-yl. Most preferably, A is benzyl.\nIn still other preferred embodiments, B is C(xe2x95x90O), OC(xe2x95x90O), or S(xe2x95x90O)2. In more preferred embodiments, B is C(xe2x95x90O) or OC(xe2x95x90O). In especially preferred embodiments, A is benzyl and B is OC(xe2x95x90O).\nIn further preferred embodiments, R3 is H, n is 0, 1, or 2, and M is sodium. Most preferably, n is 0.\nIn still further preferred embodiments, Aaa is independently Ala, Phe, or Leu.\nAs used herein, the term xe2x80x9calkylxe2x80x9d includes straight-chain, branched and cyclic hydrocarbon groups such as, for example, ethyl, isopropyl and cyclopentyl groups. Preferred alkyl groups have 1 to about 10 carbon atoms. xe2x80x9cCycloalkylxe2x80x9d groups are cyclic alkyl groups. xe2x80x9cArylxe2x80x9d groups are aromatic cyclic compounds including but not limited to phenyl, tolyl, naphthyl, anthracyl, phenanthryl, pyrenyl, and xylyl. Preferred aryl groups include phenyl and naphthyl. The term xe2x80x9ccarbocyclicxe2x80x9d, as used herein, refers to cyclic groups in which the ring portion is composed solely of carbon atoms. The term xe2x80x9clower alkylxe2x80x9d refers to alkyl groups of 1 to about 6 carbon atoms. The term xe2x80x9chalogenxe2x80x9d refers to F, Cl, Br, and I atoms. The term xe2x80x9caralkylxe2x80x9d denotes alkyl groups which bear aryl groups, for example, benzyl groups. As used herein, xe2x80x9calkoxyxe2x80x9d groups-are alkyl groups linked through an oxygen atom. Examples of alkoxy groups include methoxy (xe2x80x94OCH3) and ethoxy (xe2x80x94OCH2CH3) groups. In general, the term xe2x80x9coxyxe2x80x9d when used as a suffix denotes attachment through an oxygen atom. Thus, alkoxycarbonyl groups are carbonyl groups which contain an alkoxy substituent, i.e., groups of general formula xe2x80x94C(xe2x95x90O)xe2x80x94Oxe2x80x94R, where R is alkyl. The term xe2x80x9caryloxyxe2x80x9d denotes an aryl group linked through an oxygen atom, and the term xe2x80x9caralkyloxyxe2x80x9d denotes an aralkyl group linked through an oxygen atom.\nThe terms xe2x80x9cheterocyclexe2x80x9d, xe2x80x9cheterocyclylxe2x80x9d, and xe2x80x9cheterocyclicxe2x80x9d refer to cyclic groups in which a ring portion includes at least one heteroatom such as O, N or S. Heterocyclic groups include xe2x80x9cheteroarylxe2x80x9d as well as xe2x80x9cheteroalkylxe2x80x9d groups. xe2x80x9cHeterocycloalkylxe2x80x9d denotes a heterocycle attached through a lower alkyl group. The term xe2x80x9cheteroarylxe2x80x9d denotes aryl groups having one or more hetero atoms contained within an aromatic ring. The term xe2x80x9cheteroarylalkylxe2x80x9d denotes a heteroaryl group attached through an alkyl group. The term xe2x80x9cheteroalkylxe2x80x9d denotes a heterocyclic group which contains at least one saturated carbon atom in a heterocyclic ring. Examples of heteroalkyl groups include piperidine, dihydropyridine, and tetrahydroisoquinolin-2-yl groups.\nAs used herein, the term xe2x80x9camino acidxe2x80x9d denotes a molecule containing both an amino group and a carboxyl group. As used herein the term xe2x80x9cL-amino acidxe2x80x9d denotes an xcex1-amino acid having the L-configuration around the xcex1-carbon, that is, a carboxylic acid of general formula CH(COOH)(NH2)xe2x88x92(side chain), having the L-configuration. The term xe2x80x9cD-amino acidxe2x80x9d similarly denotes a carboxylic acid of general formula CH(COOH) (NH2)xe2x88x92(side chain), having the D-configuration around the xcex1-carbon. Side chains of L-amino acids include naturally occurring and non-naturally occurring moieties. Nonnaturally occurring (i.e., unnatural) amino acid side chains are moieties that are used in place of naturally occurring amino acid sidechains in, for example, amino acid analogs. See, for example, Lehninger, Biochemistry, Second Edition, Worth publishers, Inc, 1975, pages 73-75. One representative amino acid side chain is the lysyl side chain, xe2x80x94(CH2)4xe2x80x94NH2. Other representative xcex1-amino acid side chains are shown below in Table 1.\nThe amino acid substituents denoted xe2x80x9cAaaxe2x80x9d in the compounds of Formula I can be identical to each other, or can be different from each other.\nCompounds of Formula I contain a moiety of formula: \nwherein the symbol xe2x80x9coxe2x80x9d denotes a carbon atom which can be in either stereochemical configuration (i.e., R or S). Thus, included within the present invention are compounds of Formula I wherein the carbon atom denoted by xe2x80x9coxe2x80x9d has the R configuration, compounds of Formula I wherein the carbon atom denoted by xe2x80x9coxe2x80x9d has the S configuration, and any mixtures thereof.\nFunctional groups present in the compounds of Formula I may contain blocking groups. Blocking groups are known per se as chemical functional groups that can be selectively appended to functionalities, such as hydroxyl groups, amino groups, thio groups, and carboxyl groups. Protecting groups are blocking groups which can be readily removed from functionalities. These groups are present in a chemical compound to render such functionality inert to chemical reaction conditions to which the compound is exposed. Any of a variety of protecting groups may be employed with the present invention. Examples of such protecting groups are the benzyloxycarbonyl (Cbz; Z), t-butoxycarbonyl, methyl ester, and benzyl ether groups. Other preferred protecting groups according to the invention may be found in Greene, T. W. and Wuts, P. G. M., xe2x80x9cProtective Groups in Organic Synthesisxe2x80x9d 2d. Ed., Wiley and Sons, 1991.\nFurther blocking groups useful in the compounds of the present invention include those that bear acyl, aroyl, lower alkyl, alkanesulfonyl, aralkanesulfonyl, or arylsulfonyl substituents on their amino groups.\nThe peptidyl-2-amino-1-hydroxyalkanesulfonic acids of the invention have the unexpected advantage of excellent solubility in aqueous media, preferably approximately between 1 and 300 mg/mL, and preferably at temperatures between approximately 0xc2x0 and 40xc2x0 C., with 25xc2x0 C. being especially preferred. The compounds of the invention are preferably soluble at concentrations greater than about 5 mg/mL. More preferably, the compounds of the invention are soluble at concentrations above about 20 mg/mL. Thus, the compounds can be easily administered parenterally, for example intravenously, by bolus or infusion in the clinical or laboratory setting in media such as optionally buffered water, aqueous saline, or aqueous saline buffers at approximately pH 5 to 8, preferably pH 6.5 to 7.5. Examples of these media include water, optionally buffered 0.85% (145 mM) sodium chloride, phosphate buffered saline (10 mM phosphate, 120 mM NaCl, 2.7 mM KCl, pH 7.4), and Locke\"\"s buffer (10 mM HEPES, 154 mM NaCl, 5.6 mM KCl, 2.3 mM CaCl2, 5 uM glycine, 20 mM glucose, 1 mM sodium pyruvate, pH 7.4). The enhanced aqueous solubility of the compounds of the invention is a significant improvement over the solubilities of many other cysteine protease inhibitors. For example, the aqueous solubility of the analogous dipeptide aldehyde and xcex1-ketoamide inhibitors are typically less than 1 mg/mL. See Harbeson et al. J. Med. Chem. 1994, 37, 2918-2929.\nBecause the peptidyl-2-amino-1-hydroxyalkanesulfonic acids of the invention inhibit cysteine proteases, they can be used in both research and therapeutic settings. Inhibition of cysteine protease activity can be measured using compounds of the invention. Thus, in a research environment, preferred compounds having defined attributes can be used to screen for natural and synthetic compounds which evidence similar characteristics in inhibiting protease activity. The compounds of the invention also can be used in the refinement of in vitro and in vivo models for determining the effects of inhibition of particular proteases on particular cell types or biological conditions. In a therapeutic setting, given the connection between cysteine proteases and certain defined disorders, compounds of the invention can be utilized to alleviate, mediate, reduce and/or prevent disorders which are associated with abnormal and/or aberrant activity of cysteine proteases.\nIn preferred embodiments, compositions are provided for inhibiting a cysteine protease comprising a compound of the invention. In other preferred embodiments, methods are provided for inhibiting cysteine proteases comprising contacting a protease selected from the group consisting of cysteine proteases with an inhibitory amount of a compound of the invention.\nThe disclosed compounds of the invention are useful for the inhibition of cysteine proteases. As used herein, the terms xe2x80x9cinhibitxe2x80x9d and xe2x80x9cinhibitionxe2x80x9d mean having an adverse effect on enzymatic activity. An inhibitory amount is an amount of a compound of the invention effective to inhibit a cysteine protease. The term xe2x80x9creversiblexe2x80x9d, when used to modify xe2x80x9cinhibitxe2x80x9d and xe2x80x9cinhibitionxe2x80x9d, means that such adverse effect on catalytic activity can be reversed. The term xe2x80x9cirreversiblexe2x80x9d, when used to modify xe2x80x9cinhibitxe2x80x9d and xe2x80x9cinhibitionxe2x80x9d, means that such adverse effect on catalytic activity cannot be reversed.\nCompounds provided herein can be formulated into pharmaceutical compositions by admixture with pharmaceutically acceptable nontoxic excipients and carriers. As noted above, such compositions may be prepared for use in parenteral administration, particularly in the form of liquid solutions or suspensions; or oral administration, particularly in the form of tablets or capsules; or intranasally, particularly in the form of powders, nasal drops, or aerosols; or dermally, via, for example, trans-dermal patches; or prepared in other suitable fashions for these and other forms of administration as will be apparent to those skilled in the art.\nThe composition may conveniently be administered in unit dosage form and may be prepared by any of the methods well known in the pharmaceutical art, for example, as described in Remington\"\"s Pharmaceutical Sciences (Mack Pub. Co., Easton, Pa., 1980). Formulations for parenteral administration may contain as common excipients sterile water or saline, polyalkylene glycols such as polyethylene glycol, oils and vegetable origin, hydrogenated naphthalenes and the like. In particular, biocompatible, biodegradable lactide polymer, lactide/glycolide copolymer, or polyoxyethylene-polyoxypropylene copolymers may be useful excipients to control the release of the active compounds. Other potentially useful parenteral delivery systems for these active compounds include ethylene-vinyl acetate copolymer particles, osmotic pumps, implantable infusion systems, and liposomes. Formulations for inhalation administration contain as excipients, for example, lactose, or may be aqueous solutions containing, for example, polyoxyethylene-9-lauryl ether, glycocholate and deoxycholate, or oily solutions for administration in the form of nasal drops, or as a gel to be applied intranasally. Formulations for parenteral administration may also include glycocholate for buccal administration, a salicylate for rectal administration, or citric acid for vaginal administration. Formulations for transdermal patches are preferably lipophilic emulsions.\nThe materials of this invention can be employed as the sole active agent in a pharmaceutical or can be used in combination with other active ingredients. The concentrations of the compounds described herein in a therapeutic composition will vary depending upon a number of factors, including the dosage of the drug to be administered, the chemical characteristics (e.g., hydrophobicity) of the compounds employed, and the route of administration. In general terms, the compounds of this invention may be provided in an aqueous physiological buffer solution containing about 0.1 to 30% w/v compound for parenteral administration. Typical dose ranges are from about 1 xcexcg/kg to about 1 g/kg of body weight per day; a preferred dose range is from about 0.01 mg/kg to 100 mg/kg of body weight per day. The preferred dosage of drug to be administered is likely to depend on such-variables as the type and extent of progression of the disease or disorder, the overall health status of the particular patient, the relative biological efficacy of the compound selected, and formulation of the compound excipient, and its route of administration.\nAs used herein, the term xe2x80x9ccontactingxe2x80x9d means directly or indirectly causing at least two moieties to come into physical association with each other. Contacting thus includes physical acts such as placing a compound of the invention together with a protease in a container, or administering a compound of the invention to a patient. Thus, for example, administering a compound of the invention to a human patient evidencing a disease or disorder associated with abnormal and/or aberrant activity of such proteases falls within the scope of the definition of the term xe2x80x9ccontactingxe2x80x9d.\nThe invention is further illustrated by way of the following examples which are intended to elucidate the invention. These examples are not intended, nor are they to be construed, as limiting the scope of the disclosure."}
{"text": "1. Field of the Invention\nThe present invention relates generally to a surgical instrument for insertion into human tissue and a method for use thereof.\n2. Description of the Prior Art\nInstruments and methods exist to penetrate and dilate tissue. One problem with conventional instruments is that when the instrument includes multiple components, it is often difficult to efficiently and reliably secure the components to one another. Another problem is that when multiple components have multiple handles, the disengagement of one component from another is more cumbersome, resulting in surgical inefficiency.\nIn surgical procedures involving sequential dilation of tissue, it is typical to use multiple instruments in succession to dilate the tissue to meet a desired surgical objective. For example, a guide wire is often inserted first, followed by a larger diameter stylet, which is then removed, followed by yet another instrument having a larger diameter to increase the diameter of the opening. Successive insertion and removal of multiple instruments is a time-consuming process that can hamper the surgical procedure. Additionally, the positioning of separate smaller instruments within larger instruments may lead to imprecise results as it is more difficult to maintain the instruments coaxial with one another, or limit the depth of insertion of the instruments relative to one another.\nAccordingly, there exists a need for a surgical instrument having a configuration so that multiple components may be individually or in combination engaged to a handle for easier insertion and removal of the surgical instrument and its components from a patient. There also exists a need for a surgical instrument capable of performing sequential dilation without the need to successively insert and remove a multitude of individual instruments."}
{"text": "This invention relates to authentication of users of electronic information and transaction processing systems, and more specifically to systems and methods for authenticating users of program objects in distributed computing environments based on negotiated security contexts.\nMany computer and communication systems restrict access to authorized users. As typically shown in FIG. 1, a user 110 accesses such a system 120 through a suitable interface such as a computer 130 executing a client application. The computer 130 and client application can communicate with the system 120 by either a direct connection or via the Internet using a convenient protocol such as http as illustrated by connection 140. In accessing the system 120, the client application recognizes that a user authentication process must be carried out as a threshold step, and thus the client application usually requests the user 110 to enter a logon ID and a password that uniquely identify the user 110 to the system 120. The logon ID and password are conventionally forwarded to a logon component 122 via an application server component 124 included in the system 120.\nThe logon component 122 compares the logon ID and password received from the user to an archive of logon IDs and passwords stored in a rules database 126. Upon finding a match with the forwarded logon ID, the logon component 122 retrieves the corresponding password from the rules database 126 and compares the retrieved password with the password forwarded by the user 110. If the passwords match, the logon component 122 sends an instruction or a message via the application server component 124 to the client application in computer 130 that indicates that the user 110 been properly identified and authenticated to proceed. This authentication step and/or subsequent instructions or messages may initiate a secure communication session using a convenient protocol such as https indicated by connection 150. Subsequent communication between the system 120 and computer 130 can then proceed in private.\nWhere encryption is employed, a client cryptographic token such as an electronic circuit card conforming to the standards of the PC Memory Card Interface Association (a PCMCIA card or PC Card) can be used in the computer 130. In general, a token is a portable transfer device that is used for transporting keys, or parts of keys. It will be understood that PC Cards are just one form of delivery mechanism; other kinds of tokens may also be used, such as those conforming to RSA Laboratories' software token Public-Key Cryptography Standard (PKCS) #12, floppy diskettes, and Smart Cards.\nIf the logon ID provided by the user 110 does not match an ID in the rules database 126 or if the password comparison fails, the logon component 122 typically sends a message or instruction through the application server component 124 to the client application to inform the user 110 that the submitted logon information was incorrect and to prompt the user to re-enter it. This process of entering and attempting to verify the logon information may be permitted to occur a few times, but in the event of repeated failure, the logon component 122 may finally reject further logon attempts by the user 110, direct the client application to inform the user 110 that the logon process has failed, terminate the communication session, and lock out the user from any further logon attempts.\nA password is one form of identification that may be presented to the logon component 122 that authenticates the user's access rights, which may range from simply viewing selected records in the system 120 to performing all transactions permitted by the system 120. This kind of secured transaction processing is typically “state-full” in that it maintains, in the transaction session, the process state and content of the user's logon access information. Different transactions are typically implemented in modern distributed, nested, transaction processing systems by different program objects, e.g., applications or subroutines, that are substantially independent, even to the extent of executing on different processor hardware. For a user to migrate from one secured transaction to another, which is often necessary for even simple uses of today's systems, the user is generally required to logon (i.e., be authenticated) to each transaction, often with ID's and passwords unique to each transaction. This is because state-full systems impose state routing restrictions on users, and only sessions with particular restrictions can service a given user without having to close one program object and open another, with the corresponding requisite logon.\nBesides the burden on system resources imposed by each logon, which requires access to and processing by a logon component and a rules database, state-full systems often compel each user to close one secured transaction (program object or application) before entering another, limiting the flexibility of the system from the user's perspective. In addition, if the path to the rules database is closed or if excessive traffic slows processing or access to any of the necessary components of the system, the user access to the desired application is compromised, even if the user's access is fully authorized. This becomes a significant problem for systems having many potential users because economics often limits the system resources that can be made available.\nAnother problem is that conventional enrollment systems can be viewed as “open doorways” into an otherwise protected application in that a successful logon provides a user full access to the application and a failed logon “slams the door” on access to the application. No middle ground is generally provided, whereby a properly identified user is provided partial access to an application or transaction.\nYet another problem with systems like that depicted in FIG. 1 is the vulnerability of such systems to a hacker's or a pirate's intercepting a user's logon information at any of several points and then gaining unauthorized access to a supposedly secure system, such as an online brokerage system. One countermeasure to such interception is the application of cryptography to the data being transmitted. Public-Key Cryptography (PKC), or asymmetric cryptography, is a form of data encryption that uses a pair of cryptographic keys, each pair having a public key that is used for encryption and a private (secret) key used for decryption. Exemplary PKC algorithms, which comply with contemporary government or commercial standards, are the Digital Signature Algorithm and the Rivest-Shamir-Adleman (RSA) algorithm. The alternative to PKC is a symmetric key cryptographic system that uses the same key for encryption and decryption. Exemplary symmetric systems are the Data Encryption Standard (DES) and its improvement, the Advanced Encryption Standard (AES), recently announced by the National Institute of Standards and Technology (NIST). Symmetric key cryptography is normally employed for encrypting large amounts of data since it is much faster than PKC, but PKC is still advantageously used for key distribution. Nevertheless, encrypting transmitted data may address privacy concerns in electronic commerce and communication, but encryption alone does not address the issues of integrity and authentication of the transmitted information.\nIn this application, “privacy” means the protection of a record from unauthorized access. “Integrity” means the ability to detect any alteration of the contents of a record or of the relative authority of a user to perform a transaction or access a record. “Authentication” means verification of the authority of a user to perform a transaction, use a system resource, or access an electronic record. It will be appreciated that “electronic record” and “record” mean information in any electronic form, regardless of type of medium or type of information. Thus, a record can be a tape cartridge, a voice transmission or recording, a video image, a multi-media object, a contract, metadata, a database of information, etc.\nIntegrity and authentication of information are typically handled by other cryptographic operations, in particular hashing the information to be protected and appending one or more digital signatures. In general, a one-way cryptographic function operates on the information and produces a “hash” or “message digest” in a way such that any change of the information produces a changed message digest. Since a different message digest is produced if even one bit of the information object is changed, the hash operation yields a strong integrity check. Known hashing algorithms are the Secure Hash Algorithm (SHA-1) and the Message Digest 5 (MD-5) algorithm, and new algorithms appear from time to time. Information is typically digitally signed by hashing the information, encrypting the resulting hash using the signer's private key, and appending the encrypted hash to the information. Thus, digital signatures are generated in a manner like PKC, but the keys are reversed: the encryption key is private and the decryption key is public; the digital signer signs information with the private key and a user can read the digital signature with the signer's public key. Since a digital signature is an non-forgeable data element attached or allocated to information that ties the signer to the information, the digital signature yields an authentication check. It will be appreciated that a digital signature differs from a holographic, or handwritten, signature and from a digitized holographic signature, which is a handwritten signature that has been captured electronically.\nThe uses of digital signatures typically involve uses of authentication certificates, which are non-forgeable, digitally signed data elements that bind the signers' identity information to the signers' public-key information. Authentication certificates have been standardized by the International Telecommunications Union (ITU) under International Standard X.509, as documented in “The Directory-Authentication Framework” (1988) and as interpreted by the Internet Engineering Task Force Public Key Infrastructure X.509 recommendations. An authentication certificate is digitally signed and issued by a Certification Authority that is responsible for ensuring the unique identification of all users. Each authentication certificate typically includes the following critical information needed in the signing and verification processes: a certificate version number, a serial number, identification of the Certification Authority that issued the certificate, identifications of the issuer's hash and digital signature algorithms, a validity period, a unique identification of the user who owns the certificate, and the user's public cryptographic signature verification key. A signer's authentication certificate may be appended to information to be protected with the user's digital signature so that it is possible for others to verify the digital signature.\nSingle-logon methods have been implemented in which a logon component returns a “cookie” or token to a client application that allows the client application system-wide logon in a distributed computing environment. One example of this is the SITEMINDER software product made by Netegrity, Inc., Waltham, Mass., and described at www.netegrity.com. Such single-logon methods avoid the need for repeated logons, but have severe limitations when used with state-less computing environment components.\nU.S. Pat. No. 5,757,920 for “Logon Certification” and U.S. Pat. No. 5,999,711 for “Method and System for Providing Certificates Holding Authentication and Authorization Information for Users/Machines”, both to Misra et al., describe logon certificates that are provided to support disconnected operation in distributed computing systems. Each logon certificate is a secure package holding credentials information sufficient to establish the identity and rights and privileges for a user or a machine in a domain that is not the user's/machine's home domain.\nU.S. Pat. No. 5,241,594 to Kung for “One-Time Logon Means and Methods for Distributed Computing Systems” describes storing password files in all networked computers in a distributed system and, after a user logs on to a computer, forwarding authentication information to a second computer using a secure transport layer protocol if the user wishes to use services at the second computer. The second computer compares the user's authentication information it receives with the user's authentication information it stores, and if the informations match, the user is logged on to the second computer.\nOther logon methods and systems are described in U.S. Pat. No. 5,655,077 to Jones et al. for “Method and System for Authenticating Access to Heterogeneous Computing Services”; U.S. Pat. No. 5,689,638 to Sadovsky for “Method for Providing Access to Independent Network Resources by Establishing Connection Using an Application Programming Interface Function Call Without Prompting the User for Authentication Data”; U.S. Pat. No. 5,768,504 to Kells et al. for “Method and Apparatus for a System Wide Logan [sic] in a Distributed Computing Environment”; U.S. Pat. No. 5,774,650 to Chapman et al. for “Control of Access to a Networked System”; U.S. Pat. No. 5,884,312 to Dustan et al. for “System and Method for Securely Accessing Information from Disparate Data Sources through a Network”; and U.S. Pat. No. 6,178,511 to Cohen et al. for “Coordinating User Target Logons in a Single Sign-On (SSO) Environment”.\nThe problems with systems like that shown in FIG. 1 are keenly felt in many computer and communication systems, including as just one example those employed in electronic commerce. As paper documents that have traditionally recorded transactions, such as the purchase of an object, the withdrawal of bank funds, or the execution of a contract, are replaced by electronic records, serious issues of physical control of the electronic records and access to them are raised. Systems and methods for providing a verifiable chain of evidence and security for the transfer and retrieval of electronic records and other information objects in digital formats have been described in U.S. Pat. No. 5,615,268; U.S. Pat. No. 5,748,738; and U.S. Pat. No. 6,237,096; all to Bisbee et al., and U.S. patent application Ser. No. 09/452,928, filed on Dec. 2, 1999, and Ser. No. 09/737,325, filed on Dec. 14, 2000, both by Bisbee et al. These patents and applications are expressly incorporated here by reference, and describe among other things flexible business rules that enable users to have roles that are required or enabled only at particular points in a transaction or process. For example, a user may have a role of title agent only after a transaction has closed.\nSuch work flows and processes can be more complex than those typically associated with single-logon techniques. Moreover, many electronic records available to online inquiry are neither encrypted, nor hashed, nor digitally signed since to do so increases the processing time and resources needed for authorized users to access such information."}
{"text": "The present invention relates to an apparatus for the storage and release of information regarding an article or the like, this apparatus comprising a computer, a storage unit for storing the information, at least one output unit and an interrogator unit.\nEuropean Patent Application No. 00 99 072 describes an apparatus for checking the departure or arrival of goods. By means of a hand-held reader and an identification means, an identification which accompanies or characterized each article can be scanned or interrogated. A keyboard provided in a portable terminal provides additional interrogation capability. Any information of interest may be retrieved by actuating one of the keys of the keyboard. A signal is transmitted through a wireless transmission system from the terminal to a relay station, and the signals received are passed thence, through wiring, to an exchange. The terminal comprises an alphanumeric indicator panel to display the type of product, price and remaining stock. The operation of this apparatus requires adequate knowledge and it is intended for use only by the staff of a department store or the like. A stranger, for example a customer in a department store, would not be able to use it. Moreover, such an apparatus can provide the user with only a small amount of information.\nGerman Offenlegungsschrift No. 30 30 923 describes a device designed to produce a control tape for machine tools. This device contains a digitalizing desk upon which a drawing may be placed and may be scanned with a magnifying glass and a scanning pen. The geometrical data thus digitalized may be displayed on a graphical output unit. Furthermore, corrections, amendments, (etc . . . ) may be carried out with an interactive procedure. A device of this kind is not suitable for transmitting information to a customer in a department store.\nBritish Patent Application No. 20 80 990 discloses an apparatus comprising a display screen and a tactile screen allowing a customer to carry on a direct data dialogue with a computer and a storage unit. The tactile screen, comprising a series of transparent, capacitive-action contact paths or sensors, is arranged in front of the display screen. By means of the sensors, a user may make a selection from a menu offered on the monitor and, by means of a subsequent microprocessor, another menu or the like may be retrieved from the storage unit and may be presented on the display screen. With a tactile screen of this kind, problems with accuracy and resolution arise because of the relatively small dimensions of the picture tube and the display screen. Furthermore, a user may very easily miss the desired position with his finger and thus obtain incorrect information.\nIn department stores and supermarkets, (etc . . . )there is a considerable need to provide the customer with information regarding products which are offered, the location of such products, their use, and so on. The information carried on the packages of products is often inadequate in practice, especially since, in order to reduce environmental pollution, such packages are ever becoming smaller or are being eliminated altogether. Information and advice provided directly by sales counsellors is costly and requires additional staff. Moreover in exhibitions, museums, airports, etc . . . there is also a considerable need for information regarding objects on display, the location of such objects, and flight connections or the like."}
{"text": "The following relates to the statistical sampling arts, Monte Carlo sampling arts, Markov Chain Monte Carlo sampling arts, computational design and analysis arts, and so forth, and to arts employing same such as machine translation, natural language processing, and so forth.\nStatistical or Monte Carlo sampling employs random sampling of a domain to estimate a distribution, to obtain samples in (approximate) accordance with a distribution, to compute an integral, or to generate some other result of interest. In Markov Chain Monte Carlo (MCMC) approaches, a Markov chain of samples is constructed over a chain length sufficient to converge to an equilibrium result. For example, in a “random walk” MCMC approach, the chain is constructed by moving in a random direction from a currently accepted sample to a next sample. The result generated by a random walk MCMC approach can be adversely impacted by the choice of starting point for the random walk, which can lead to slow convergence or to convergence to a local extremum.\nIn Metropolis-Hasting (MH) MCMC approaches, the random walk is replaced by sampling in accordance with a normalized “proposal” distribution that is chosen such that the domain is readily sampled in accordance with the proposal distribution. In MH sampling, the Markov chain is constructed as follows: given a currently accepted sample, a new sample is chosen in accordance with the proposal distribution. If the new sample satisfies an acceptance criterion then it becomes the next accepted sample in the chain; otherwise, the currently accepted sample is repeated in the chain. In probabilistic sampling applications, the acceptance criterion is suitably embodied as a target distribution (not necessarily normalized) and the acceptance or rejection is performed on a statistical basis in accordance with the target distribution, for example by comparing the target distribution value at the node with a value drawn from a uniform probability. The MH approach eliminates dependence on the starting point of the sampling and, for a suitable choice of the proposal distribution, can ensure convergence to the globally optimal result in the limit. However, a poor choice for the proposal distribution can still lead to slow convergence, and there may be no suitable basis for choosing a “good” proposal distribution that ensures suitably fast convergence.\nFor example, consider MH sampling guided by an unnormalized target distribution that determines whether a sample is accepted or rejected. (Said another way, the acceptance criterion is statistical in nature and is defined by the unnormalized target distribution). Here, the convergence rate is typically controlled by the similarity (or dissimilarity) of the proposal distribution to the target distribution, as this determines how frequently samples generated by the proposal distribution are accepted by the target distribution. If the proposal distribution differs significantly from the target distribution then the MH sampling may become “trapped” for extended periods at a currently accepted sample, as the proposal distribution has difficulty proposing a new sample that is “acceptable” to the target distribution. In practice, however, it can be difficult to design a normalized proposal distribution that is sufficiently similar to the unnormalized target distribution to provide fast convergence.\nIn independent MH sampling, the proposal distribution is independent of the last accepted sample. In generalized MH sampling, the proposal distribution is conditioned on the last accepted sample. The generalized MH sampling is thus in a sense “closer” to the random walk MCMC approach, insofar as the next sample has some dependency on the last accepted sample. Accordingly, generalized MH sampling is also sometimes referred to as random walk MH sampling. The generalized MH sampling approach can improve convergence since the conditioning upon the last accepted sample tends to increase similarity between the conditional proposal distribution and the unnormalized target distribution. However, as with random walk MCMC, conditioning the proposal distribution on the last accepted sample can lead to undesirable dependence of the result on the starting point. Moreover, even with conditioning on the last accepted sample, the proposal distribution may still have substantial dissimilarity from the target distribution leading to slow convergence of the generalized MH sampling.\nThe following sets forth improved methods and apparatuses."}
{"text": "Videoconferencing between multiple people typically involves many pieces of hardware working together, such as cameras, microphones, headsets, computers, monitors, routers, and other equipment. As the number of different devices increases, the difficulty and time consumption in identifying the source of a problem also increases, as the problem may lie in any one of the devices, and may also include network limitations. A problem at the beginning of a videoconference results in inefficiency and annoyance among conference participants, as well as other personnel who may be responsible for setting up and operating the videoconference equipment.\nCurrent solutions typically involve a videoconference equipment manufacturer providing a specific hardware device to connect to the manufacturer's videoconference device to test for video integrity, such as detecting packet loss and jitters of the videoconference stream. However, these manufacturer's solutions are typically limited to the manufacturer's dedicated hardware and limited by the type of information provided by the manufacturer's dedicated hardware.\nFor the foregoing reasons there is a need for a platform agnostic monitoring solution that can detect the health of a system of devices and each individual device so problems may be readily identified and targeted."}
{"text": "Teaching people to read, especially children, fosters greater literacy and knowledge among the general populous. Getting children to read, unfortunately, can be a laborious task. Many reading materials are unattractive to the eye and do not foster imagination among the potential reader. The less than appealing designs of reading materials often result in an apathetic reader or avoidance of reading altogether.\nTo encourage reading, some reading materials have attractive and/or fanciful color designs or pictures to entice the reader. Often, however, these materials are located only on the exterior cover of the materials causing the reader to lose interest on the interior pages. Reading materials with pictures or design work on the interior pages are common but have inherent problems in themselves. These problems include expensive printing costs, and thick volume size due to the increased number of pages due to the artwork.\nCommon reading materials also have another significant drawback. For the material to be seen, a proper viewing perspective must be maintained by the reader (i.e. the book must be in front of the reader.) A young child must often crowd next to an adult to view the text and/or artwork in the book. If multiple children are present, viewing can become even more restrained as each child tries to obtain a proper perspective of the book. The reader must often stop and show the pages to the children, slowing the pace of the story, and affecting interest in the work for all involved. In large families or in classroom situations, the problems presented are self-evident.\nChildren retelling a previously read story is also a significant part of the reading experience. The retelling of the story by the child fosters memory recall and imagination. If a child does not have visual cues upon which to initiate successive steps in the story, the child may not want or be able to retell the story. Reading materials without visual prompts, therefore, lack a critical feature in the learning process.\nIt is therefore desired to have reading materials which do not have these drawbacks and which encourage reading, especially among younger people."}
{"text": "Distilling units of this kind are used for the most widely varying analyses in laboratories. The system, being closed during the distilling operation, is particularly important when substances which are gaseous or which evaporate quickly are to be bonded or combined by the distilling receiver. The invention will be described hereinafter with reference to a distilling unit for determining nitrogen in samples of milk, by direct distillation. However, apparatus according to the invention can be used for any distillation unit of the kind indicated.\nDetermining nitrogen by direct distillation is based on the fact that a sample, upon being boiled in alkaline solution, gives off ammonia. The major proportion of the ammonia originates from the rapid hydrolysis, which is concluded after a few minutes, of the albumin components, glutamine and asparagine. This fact makes it possible to determine experimentally the relationship between total nitrogen or protein respectively and the ammonia nitrogen which is liberated upon boiling in an alkaline solution. The whole analysis operation is reduced, by virtue of such direct distillation, to the following steps: introduction by means of pipette, dilution, alkaline ajustment, distillation, titration and calculation. It is important however that the distillation time is strictly observed and that the distillate discharge pipe is immersed into the receiver so that even gaseous ammonia components are combined in the distillate receiver.\nWhen the distilling operation is stopped, there occurs a reduced pressure which, when the distillate discharge pipe is immersed, has the result that the distillate is sucked back into the distilling vessel. In order to avoid this, the receiver with the distillate must be removed precisely after the predetermined distilling time, before distillation is interrupted.\nEquipment is already known for substantially automating the distillation operation, in order to achieve the most rational sample processing possible. With such equipment, after the preset distillation time has expired, the distilling receiver is automatically lowered before the distillation operation is interrupted. This mode of operation reliably prevents the distillate from being drawn back into the distilling vessel. However, this equipment suffers from the disadvantage that the lifting means and the electronic control arrangement which is required for that purpose is costly and complicated. The lifting mechanism is also liable to trouble and requires careful maintenance and care."}
{"text": "1. Field of the Invention\nThe present invention relates to orthopedics, and, particularly, to orthopedic implants.\n2. Description of the Related Art\nDuring orthopedic surgery, such as a total hip arthroplasty, orthopedic prostheses are implanted into a patient's body. For example, a femoral stem may be implanted into the intramedullary canal of the patient's femur. After the stem is inserted into the intramedullary canal at a desired depth and orientation, a femoral head may be attached to the proximal end of the stem. The femoral head fits into a substantially hemispherically shaped socket of an acetabular prosthesis and provides a surface for articulation between the femoral head and the acetabular prosthesis.\nIn certain embodiments, acetabular prostheses comprise two separate components, an acetabular backing component and an acetabular articulating component, which are often referred to as a shell and liner, respectively. The backing component is generally hemispherically shaped and is affixed and embedded into an acetabulum of a patient. Similarly, the liner has a hemispherical shape configured to mate with an internal cavity of the backing component. The liner further includes a socket to mate with a femoral head. The acetabular backing components are typically formed of solid metal, such as titanium or a titanium alloy, or from stainless steel. In contrast, the acetabular articulating components may be formed from an ultra-high molecular weight polyethylene polymer material (“UHMWPE”). These types of acetabular cups are often referred to in the art, and sometimes herein, as “metal backed” components.\nMetal backed orthopedic components generally experience contact stresses during normal joint articulation. For example, in conventional metal backed acetabular components, previous studies have shown that polymer thickness and conformity between the femoral head and the polyethylene liner play a significant role in the level of contact stress experienced. For highly conforming acetabular component designs, the contact stress in the polymer liner is very sensitive to the polymer thickness when the polymer thickness is small. Specifically, it has been demonstrated that when the polymer thickness is less than 4-6 mm, the contact stress generally increases rapidly as the thickness decreases. Therefore, a minimum polymer thickness of 4-6 mm has been maintained for metal backed acetabular components in an attempt to lessen the contact stress. Further, having a 4 mm minimum thickness of the polyethylene liner in an acetabular component has been widely accepted as an industry design rule and has been traditionally relied upon in the manufacture of acetabular components.\nStress shielding is a phenomenon that can potentially occur with metal backed orthopedic components and can occur when a portion of the stress normally exerted on a patient's bone is instead borne by the orthopedic component, such as an acetabular cup. As result, the bone may begin to undergo atrophy and decalcification resulting in a weaker bone structure. In severe stress shielding, the bone may be resorbed by the body, decreasing the amount of bone stock at the fracture site."}
{"text": "This invention pertains to a leveling tool and, more particularly, to a tool for providing a smooth, inclined surface about a drain or the like.\nIn concrete/mud work it is desirable to provide an inclined, finished surface about a drainpipe or the like, the surface being inclined about the drainpipe so as to direct water thereto. Such surface may either then remained unfinished, as in an outside drain, or have tile mounted thereon as in a bathroom. It is thus desirable to provide a leveling tool which allows the worker to easily provide this inclined, smoothed surface.\nIn response thereto we have provided a leveling tool in the form of an elongated beam having a centrally located seat which receives the open end of-the drainpipe therein. Extending from this central seat in opposed directions is a pair of inclined surfaces. Each surface has an angle corresponding to the desired angle of inclination of the smoothed surface about the drainpipe. After the cement or other xe2x80x9cmudxe2x80x9d-like material has been poured about the drainpipe, the drainpipe is seated in the central seat such that the inclined surfaces extend away from the drainpipe in contact with the unfinished mud. User rotation of the device about the drainpipe, which acts as an axis of rotation, smoothes the poured material surface therebelow. During rotation of the device, pockets in the sidewall pick up the excess mud for transport away from the smoothed surface. Bubble levels in the device indicate that the device is at a proper position during rotation. A plate can be positioned within the seat so as to regulate the depth of penetration of the drainpipe within the seat and thus the depth of the smoothed surface below the top edge of the drainpipe.\nIt is therefor a general object of the invention to provide a device with leveling surfaces for presenting a smoothed surface about a drainpipe or the like at a desired inclination relative thereto.\nAnother object of this invention is to provide a device, as aforesaid, which presents a seat releasably engageable with the open end of the drainpipe, the leveling surfaces being rotatable about the drainpipe.\nA still further object of this invention is to provide a device, as aforesaid, wherein the degree of penetration of the drainpipe into the seat can be adjusted.\nAnother object of this invention is to provide a device, as aforesaid, wherein pockets in the sidewall of the device pick up the excess mud material therein during smoothing of the poured material surface.\nA particular object of the invention is to provide a device, as aforesaid, which has bubble levels therein for properly positioning the device during operation.\nOther objects and advantages of this invention will become apparent from the following description taken in connection with the accompanying drawings, wherein is set forth by way of illustration and example, a now preferred embodiment of this invention."}
{"text": "1. Field of the Invention\nThe present invention relates to a method and to a circuit arrangement for determining the quality of virtual connections extending via input lines and output lines of an asynchronous transfer mode (ATM) switching equipment, whereby respective message cells are transmitted according to an asynchronous transfer mode process during the course of these virtual connections and each of these message cells thereby has a cell header available for identifying the respective virtual connection for the input line or, respectively, output line coming into consideration.\n2. Description of the Prior Art\nThe problem generally arises in telecommunication switching systems of monitoring the quality of defined connections in order to be able to provide subscribers with appropriate information, for example, given complaints. The operators of public or private switching equipment, however, also have the need for or, respectively, raise the requirement of being able to measure and monitor the quality of the switching equipment. Quality measurements of this type also enter into the maintenance concepts in order to be able to, for example, avoid predictable deteriorations in quality beyond an allowable tolerance limit.\nGiven a method and a circuit arrangement of the type initially set forth, it is therefore an object of the present invention to provide a way to determine the quality of virtual connections on the basis of the message cells transmitted during the course of such virtual connections.\nIn a method of the type set forth above, this object is achieved, according to the present invention, in an improved method for determining the quality of virtual connections extending via input lines and output lines of an ATM switching equipment in which respective message cells are transmitted according to an asynchronous transfer mode during the course of these virtual connections and each of the message cells thereby has a cell header available with which the respective virtual connection for the input line or, respectively, output line coming into consideration is identified, and is particularly characterized in that the plurality of message cells accepted and output are separately acquired in a call-associated manner in the ATM switching equipment on the basis of the cell headers of the message cells accepted via the input lines or, respectively, output to the output lines, and in that, in response to the cleardown of the respective virtual connection, the loss rate of the message cells is calculated as a measure for the quality of the respective virtual connection, being determined from the previously, separately-acquired plurality of accepted and output message cells.\nThe present invention yields the advantage that the quality of the individual virtual connections can be continuously determined on the basis of the transmitted message cells without significantly dynamically loading an ATM switching equipment for this determination.\nAccording to a particular feature of the invention, the aforementioned improved method is further particularly characterized in that a mean loss rate is formed as quality measure for the ATM switching equipment on the formation being from the identified loss rates of the virtual connections established within a defined time interval. Also, the improved method is particularly characterized in that the loss rate identified for the respective virtual connection is involved in the calculation of charges. This feature and development yields the advantage of a simple determination of the quality of an ATM switching equipment.\nThe advantage of the improved circuit arrangement set forth above is that only the message quantity forwarded by an ATM switching equipment to a subscriber equipment, in fact, enters into the calculation of charges. Message cells repeatedly transmitted within an ATM switching equipment or, respectively, misrouted message cells due, for example, to transmission malfunctions are left out of consideration for calculating charges.\nThe circuit arrangement for the implementation of the method set forth above is characterized in that a plurality of counters corresponding in number to the plurality of virtual connections extending via the input lines and the output lines is assigned to the input lines and to the output lines, in that the plurality of counters can be placed into a defined initial counter reading proceeding from a central control device during the respective course of the setup of the allocated virtual connection, in that the individual counters are individually drivable given the appearance of message cells on the allocated input line or, respectively, output line, being individually drivable according to the prescription of the cell header respectively contained in the message cells, and in that the momentary counter reading of a counter is variable by a counting unit proceeding from a defined initial counter reading, the reading being variable with each drive, and in that the momentary counter readings of the counter assigned to a virtual connection can be interrogated at least at the cleardown of the connection by a central control device for calculating the loss rate.\nThe advantage of the circuit arrangement is thereby that the registration of the message cells transmitted during the course of virtual connections occurs in peripheral equipment so that the central control device of an ATM switching equipment is largely relieved of the determination of the quality of the virtual connections.\nAccording to other features and developments of the invention, the circuit arrangement is particularly characterized in that the momentary counter readings of the counter can be additionally interrogated at fixed time intervals. Also, the arrangement is particularly characterized in that the counters are fashioned as periodically-circulating counters and in that a message signal is separately offered to the central control device before a cleardown of the respective virtual connection each time a defined counter reading is reached. The advantage of these features and developments is that the counters provided for the registration of message cells require only a relatively low counting capacity and, therefore, the required circuit expense can be kept low."}
{"text": "(1) Field of the Invention\nThe present invention relates to the use of specific odorants in hypochlorite-containing bleaching agents in order to avoid the unpleasant odor of such agents on human skin, in particular on the hands, as a result of contact with the agent.\nEffective cleaning is one of the requirements that influences consumer acceptance of such agents. This requires that agents having broad potential for spot removal be made available, which agents can, for example, remove oily and greasy stains from textiles or from hard surfaces such as floor tiles or wall tiles. In order to meet this requirement, agents containing (inter alia) hypochlorite bleaches have been developed, in which agents the hypochlorite contributes, as a strong oxidizing agent, to the chemical breakdown, destruction, and removal of the stains. A further advantage of using hypochlorite is that it acts as an effective disinfecting agent. One disadvantage that is perceived, however, is that after contact with human skin, which contact necessarily occurs when such agents are used with bare hands and cannot always be avoided even otherwise, for example when gloves are used, an odor occurs that clings to the skin and cannot always be entirely removed even after repeated washing with water.\n(2) Description of Related Art, Including Information Disclosed Under 37 C.F.R. Sections 1.97 and 1.98\nTo remedy this, it has been proposed, for example in European Patent Application EP 439 878, to use a volatile perfume composition.\nEuropean Patent Application EP 0 606 707 describes agents that contain a polymer component together with a hypochlorite compound, and that lead to a decrease in unpleasant odors of the agent itself and of surfaces cleaned therewith. This could also be attributed to the high viscosity of said agents, so that compounds that are responsible for the disadvantageous odor are possibly enclosed in vesicles.\nIn the aforesaid existing art, greater emphasis was therefore placed on avoiding disadvantageous odors of the actual agent and of the surfaces cleaned therewith, and less attention was paid to the problem of odor formation on skin surfaces that have come into contact with the agent.\nEuropean Patent Application EP 0 812 909 A1 describes the use of polycarboxylate polymer in hypochlorite-containing bleaching agents to avoid unpleasant odors that result from contact by such agents with skin surfaces.\nEuropean Patent Application EP 0 622 451 also deals with the matter of odor formation during and after the use of chlorine-containing agents, and proposes using a perfume. The difficulty encountered here is that, on the one hand, only a few fragrances are stable in the presence of the strong oxidizing ability of chlorine-based bleaching agents, and on the other hand, the odor of chlorine is very strong and is difficult to mask using a perfume. EP 0 622 451 proposes using, in such agents, 0.000002 wt % to 2 wt % bleach-stable fragrances, selected from the classes of the acetals, aldehydes, esters, alcohols, ketones, ethers, nitrites, and terpenes."}
{"text": "Precision frequency references are required in many types of electronic devices such as cellular phones and other handheld devices. Crystal oscillators are typically used to provide the precision frequency reference in these electronic devices. However, crystal oscillators have several inherent disadvantages including large bulky size, fragility, and high cost. In addition, the size and cost of crystal oscillators is related to the resonant frequency so that as the frequency increases, the size decreases, and the cost and fragility may rapidly increase. As the size of electronic devices continues to decrease, the use of crystal oscillators becomes more problematic due to the size, fragility, and cost limitations.\nSemiconductor oscillators have been a poor alternative to crystal oscillators and are generally unsuitable for use as a precision frequency reference due to excessive variation in the oscillating frequency, especially with changes in temperature."}
{"text": "(1) Field of the Invention\nThe present invention relates to a resynchronization control apparatus and a resynchronizing method suitable for use in a subscriber communication machine that can communicate in synchronization with a transmission cycle in, for example, ISDN (Integrated Services Digital Network) ping-pong transmission [TCM (Time Compression Multiplex) transmission] over a predetermined communication line such as a subscriber line.\n(2) Description of Related Art\nRecent widespread of multimedia-type services such as Internet and the like to the whole society including general home strongly demands early provision of an economical, highly reliable digital subscriber line transmission technique by which the user can enjoy such services. Enormous cost and time are required to newly lay communication lines, so that there have been proposed various methods for high-speed data communications using existing communication lines.\nFor instance, xDSL (Digital Subscriber line) attracts attention as a digital subscriber line transmission technique that uses existing telephone lines as high-speed data communication lines. The xDSL is a transmission system using existing subscriber lines, being also one of modulation-demodulation techniques. The xDSL is roughly classified into two according to whether the upstream transmission rate from the subscriber's premise (hereinafter referred to as subscriber's side) to the accommodating office (hereinafter referred to as office's side) and the downstream transmission rate from the office's side to the subscriber's side are symmetric or asymmetric.\nAs the symmetric type, there are HDSL (High-bit-rate DSL) whose upstream and downstream transmission rates are approximately 1.5 to 2.0 Mbps (megabit per second), SDSL (Single-line DSL) whose upstream and downstream transmission rates are approximately 160 k to 2.0 Mbps and the like, for example. As the asymmetric type, there is ADSL (Asymmetric DSL) that is tentatively operated in these years, extensively. The ADSL is further classified into “G.dmt” whose downstream transmission rate is approximately 6 Mbps, and “G.lite” (also called as light ADSL) whose downstream transmission rate is approximately 1.5 Mbps, both of which employ a distinctive modulation system called DMT (Discrete Multiple Tone) modulation.\nIn brief, the DMT modulation system divides the transmission frequency bandwidth into subcarriers each of about 4 kHz (in the case of “G.lite”, approximately a maximum of 128 carriers in the downstream although it depends on conditions), and modulates each of the subcarriers. The DMT modulation system has a feature to be resistible to noise at a specific frequency since, even when a certain subcarrier is unusable due to an effect of the noise at the specific frequency, it allows the communication using another subcarriers.\nHereinafter, an ADSL transmission system employing such the DMT modulation system will be described in detail.\n(1) Description of ADSL Transmission System\nFIG. 6 is a block diagram showing an example of the ADSL transmission system. The ADSL transmission system shown in FIG. 6 is configured with an ADSL machine 650 installed on the office's side 610 and an ADSL machine 660 installed on the subscriber's side 620, which are connected to each other over a metallic line (telephone line) 70. Incidentally, the ADSL machine 650 on the office's side 610 will be occasionally referred to as an office ADSL machine 650, whereas the ADSL machine 660 on the subscriber's side 620 a subscriber ADSL machine 660, hereinafter. When not particularly discriminated between the office's side 610 and the subscriber's side 620, they will be referred to merely as ADSL machines 650 and 660.\nAs shown in FIG. 6, the office ADSL machine 650 comprises, as a transmitter (transmitting block) 910, a serial to parallel buffer 10, an encoder 20, an IFFT (Inverse Fast Fourier Transformer) 30, a parallel to serial buffer 40, a D/A (digital/analog) converter 50, and a transmit bit map memory 60.\nOn the other hand, the subscriber ADSL machine 660 comprises a receiver (receiving block) 960 including an A/D converter 80, a TEQ (Time-domain Equalizer) 90, a serial to parallel buffer 100, an FFT (Fast Fourier Transformer) 110, an FEQ (Frequency-domain Equalizer) 120, a decoder 130, a parallel to serial buffer 140, a receive bit map memory 150, an AGC (Automatic Gain Controller) 160, and multipliers 170 and 180.\nFIG. 6 shows only a downstream structure from office's side to the subscriber's side. However, the office ADSL machine 650 is practically provided with a receiver having a function equivalent to that of the receiver 960 in the subscriber ADSL machine 660, whereas the subscriber ADSL machine 660 is provided with a transmitter having a function equivalent to that of the transmitter 910 in the office ADSL machine 650, whereby communication in the upstream is performed similarly to communication in the downstream, in principle. Here, the description will be made by way of the light ADSL (G.lite).\nIn the office ADSL machine 650, the transmit bit map memory 60 holds data (bit map) defining assignment of transmission bits to each (sub) carrier of a DMT signal to be generated. The serial to parallel buffer 10 stores transmit data that is serial data of one symbol duration (¼ kHz) converts the stored data into parallel data, and outputs the data. At this time, the assignment of a transmission bit number to each carrier (division of the frequency bandwidth) is performed according to the transmit bit map stored in the above transmit bit map memory 60. For example, when the number of subcarriers is i+1 where the subcarriers are C0 to Ci, a bit group bi of the above parallel data is assigned as a bit group to be transmitted on a subcarrier Ci in FIG. 6.\nThe above subcarriers sometimes include a carrier for synchronizing a timing called a pilot tone [a carrier positioning in the center of the 128 carriers is a pilot tone in the case of the downstream in the light ADSL], as shown in FIG. 7, for example. This pilot tone is used to transmit only a signal for synchronizing a timing; no data for the pilot tone exists in the transmit bit map memory 60.\nThe encoder 20 performs a predetermined modulation process such as quadrature amplitude modulation (QAM) on parallel data outputted from the serial to parallel buffer 10 for each of the above subcarriers according to the above transmit bit map. The IFFT 30 performs inverse fast Fourier transform on data (frequency-domain data) outputted from the encoder 20 to convert the data into time-domain data, thereby obtaining a DMT signal. Namely, the encoder 20 and the IFFT 30 function as a DMT modulating unit that DMT-modulates transmit data.\nThe parallel to serial buffer 40 converts the data (DMT signal) that has been subjected to the inverse fast Fourier transform in the above IFFT 30 into serial data, and adds a cyclic prefix to be described later thereto. The D/A converter 50 converts an output (serial data) of the parallel to serial buffer 40 to an analog signal using a predetermined sampling frequency (for example, 1.104 MHz). The obtained analog signal is outputted to the metallic line 70.\nIn the subscriber ADSL machine 660, the multiplier (for analog signal) 170 multiplies an analog signal received from the office ADSL machine 650 over the metallic line 70 by an arbitrary coefficient. The A/D converter 80 samples an output (analog signal) of the multiplier 170 at a predetermined sampling frequency (for example, 1.104 MHz) to convert the signal into a digital signal. The multiplier (for digital signal) 180 multiplies the digital signal from the A/D converter 80 by an arbitrary coefficient. The AGC 160 controls the coefficients to be multiplied in the above multipliers 170 and 180.\nMultiplying a signal by an arbitrary coefficient is equivalent to amplifying the signal. Namely, the AGC 160 controls the coefficients to be multiplied in the multipliers 170 and 180, thereby controlling an amplification gain of a signal received over the metallic line 70. Incidentally, the amplification gain (coefficient) of an analog signal in the multiplier 170 will be called an analog AGC value, whereas the amplification gain (coefficient) of a digital signal in the multiplier 180 will be called a digital AGC value.\nThe TEQ 90 is an equalizer in time domain configured with, for example, an FIR (Finite Impulse Response) filter, which performs a predetermined process so that intersymbol interference (ISI) with an inputted signal is placed in the cyclic prefix added in the parallel to serial buffer 40 (detailed of which will be described later). The serial to parallel buffer 100 removes the cyclic prefix from data outputted from the TEQ 90, converts the data to parallel data, and outputs the same.\nThe FFT 110 converts output data from the above serial to parallel buffer 100 to data in frequency domain in the fast Fourier transform. The FEQ 120 is an equalizer in frequency domain, which equalizes data converted into the data in frequency domain by the FFT 110 as above according to transmission characteristics (frequency characteristics) of the metallic line 70, thereby compensating effects on the amplitude and phase exerted when the data passes through the metallic line 70 for each carrier at a different frequency.\nThe decoder 130 performs a predetermined demodulating process (QAM demodulation or the like) on output data of the above FEQ 120. The parallel to serial buffer 140 converts parallel data outputted from the decoder 130 into serial data, and outputs it. The receive bit map memory 150 holds information (receive bit map) on the transmission bit number of each carrier assigned to each carrier on the transmitting side according to the transmit bit map in the transmit bit map memory 60. On the basis of this information, the demodulation process by the decoder 130 and the parallel to serial buffer 140 mentioned above is carried out.\nNext, description will be made of an operation of the above-structured ADSL transmission system.\nWhen transmit data is inputted to the office transmitter 910, the transmit data of one symbol duration (¼ kHz) is held in the serial to parallel buffer 10. The held data is divided into groups each of the transmit bit number per carrier determined beforehand according to the transmit bit map 60, and outputted to the encoder 20.\nThe encoder 20 converts the inputted bit sequences into signal points to be quadrature-amplitude-modulated, and output them to IFFT 30. The IFFT 30 performs the inverse fast Fourier transform on outputs of the encoder 20 to quadrature-amplitude-modulate each of the signal points, and outputs them to the parallel to serial buffer 40. Meanwhile, the encoder 20 and the IFFT 30 performs the DMT modulation.\nThe parallel to serial buffer 40 adds 16 samples (240 to 255 samples) out of outputs of the above IFFT 30 as the cyclic prefix to the head of the DMT symbol (details of which will be described later). The data to which the cyclic prefix has been added is sent from the parallel to serial buffer 40 to the D/A converter 50, converted to an analog signal at a sampling frequency of 1.104 MHz therein, and transmitted to the subscriber receiver 960 over the metallic line 70.\nIn the subscriber receiver 960, the analog signal received over the metallic line 70 is amplified by the multiplier 170, converted into a digital signal at 1.104 MHz by the A/D converter 80, and inputted to the multiplier 180. The multiplier 180 again amplifies the inputted digital signal, and outputs it to the TEQ 90. At this time, the AGC 160 measures a magnitude of each signal at the multipliers 170 and 180, and sets the AGC values to each signal and changes the same as needed.\nThe TEQ 90 equalizes an output of the multiplier 180 in time domain such that ISI is placed within the cyclic prefix of 16 samples, and makes the serial to parallel buffer 100 hold data of one DMT symbol. The serial to parallel buffer 100 removes the above cyclic prefix from the data of one DMT symbol inputted from the TEQ 90, then converts the remaining data into parallel signals, and outputs them to the FFT 110.\nThe FFT 110 performs fast Fourier transform on outputs of the above serial to parallel buffer 100 to convert signals in time domain into signal point data in frequency domain. The FEQ 120 compensates effects on the amplitude and phase of the converted signal point data on each carrier having a different frequency exerted when the signal passes through the metallic line 70. The decoder 130 then demodulates the data according to the receive bit map in the receive bit map memory 150 holding the same values as the transmit bit map 60. The data demodulated by the decoder 130 is held for awhile in the parallel to serial buffer 140, converted into a serial bit string, and outputted as receive data.\n(2) Description of Equalizer\n(2-1) Detailed Description of TEQ 90\nNext, description will be made of a role of the above TEQ 90.\nWhen the DMT symbol inputted to the above parallel to serial buffer 40 in FIG. 6 is in a state shown in FIG. 8(A), the parallel to serial buffer 40 copies 16 samples in the tail of the DMT symbol, and attaches them to the head of the DMT symbol, as shown in FIG. 8(B). The copied portion is the above cyclic prefix.\nThe DMT symbol added thereto the cyclic prefix as shown in FIG. 8(C) is sent to the D/A converter 50, converted to an analog signal at the sampling frequency of 1.104 MHz in the D/A converter 50, and transmitted to the subscriber ADSL machine 660 over the metallic line 70. The receive signal received over the metallic line 70 whose amplitude characteristics and phase characteristics to the frequency are not constant is distorted due to an effect of ISI, as shown in FIG. 8(D), for example.\nThe TEQ 90 described with reference to FIG. 6 performs such a process (ISI compressing process) as to place the ISI within only the cyclic prefix of 16 samples, and the serial to parallel buffer 100 removes the cyclic prefix, whereby a DMT symbol in which the effect of the ISI has been eliminated is obtained as shown in FIG. 8(F).\nThe TEQ 90 functions to eliminate the effect of ISI from the received signal using the cyclic prefix. In more detail, since the metallic line 70 has non-linear low pass filter (LPF) characteristics that deteriorate the transmission characteristics in the high frequency bandwidth, an impulse response having a constant length is generated at a discontinuous portion existing between symbols. This impulse response is overlaid on the data signal to cause deterioration of the signal.\nAccordingly, 16 bits in the tail are added to the head of the symbol to form the cyclic prefix. A portion at which the cyclic prefix is combined with the symbol becomes continuous, so that no impulse response is generated at this portion. On the other hand, a portion at which the head of the cyclic prefix is combined with the preceding symbol becomes discontinuous, so that the impulse response is generated at this portion.\nBy inputting the received data to the TEQ 90 having characteristics of a high pass filter (HPF) reverse to those of the metallic line 70, the impulse response can be placed within the cyclic prefix. The cyclic prefix so processed is removed, whereby data not affected by the impulse response can be obtained.\n(2-2) Adaptive Operation Algorithm of TEQ 90\nIn order that the TEQ 90 performs the above-mentioned process to place ISI within only the cyclic prefix of 16 samples as shown in FIG. 80(E), an adaptive operation to make the TEQ 90 have such the characteristics is required. For this purpose, the TEQ 90 comprises, as a block for the adaptive operation only, a reference signal generation block 920, a delay unit 930, a target channel block 940, and an adder 950, as shown in FIG. 9, for example. Incidentally, identical numerals in FIG. 9 identify identical or like parts in FIG. 6 described above. The structural elements of the receiving side block 960 shown in FIG. 6 other than the TEQ 90 are not shown in FIG. 9 in order to simplify the description.\nThe reference signal generation block 920 generates the same signal x(t) as a transmit signal (known received signal in the receiving side block 960) transmitted in the past from the transmitting side block 910, and outputs it. The delay unit 930 delays the reference signal x(t) generated by the reference signal generating block 920 by a predetermined time, and outputs it.\nThe target channel block 940 outputs a result b(t)*x(t), which is to be a target when the characteristics of the TEQ 90 are adjusted, obtained by convolution-integrating a characteristic b(t) (provided a delay quantity caused by the metallic line 70 is excluded) synthesized from the characteristics of the metallic line 70 and the characteristics of the TEQ 90 for the reference signal x(t) delay by the delay unit 930. The adder 950 computes a difference between an output z(t) of the TEQ 90 and an output b(t)*x(t) of the target channel block 940, and supplies the obtained result e(t) to the TEQ 90 and the target channel block 940. The TEQ 90 and the target channel block 940 thereby perform the adaptive operation such that the output e(t) of the adder 950 becomes “0.”\nNext, description will be made of the adaptive operation of the above TEQ 90.\nWhen a transmit signal x(t) is transmitted from the transmitting side block 910, the transmit signal x(t) is received by the receiving side block 960 over the metallic line 70. In the receiving side block 960, a result z(t) obtained by adding the characteristics of the TEQ 90 to the received signal by the TEQ 90 is supplied to the adder 950.\nAt this time, the reference signal generation block 920 generates a reference signal x(t) that is assumed to be the same as the transmit signal, and outputs it. The delay unit 930 delays the reference signal x(t) by a predetermined quantity such that the phase of an output z(t) of the TEQ 90 coincides with the phase of an output b(t)*x(t) of the target channel block 940. The target channel block 940 convolution-integrates the characteristic b(t) for the reference signal x(t), and supplies the obtained result b(t)*x(t) to the adder 950.\nThe adder 950 computes a difference e(t) between the output z(t) of the TEQ 90 and the output b(t)*x(t) of the target channel block 940, and supplies the obtained result to the TEQ 90 and the target channel block 940. The TEQ 90 and the target channel block 940 perform the adaptive operation on the basis of the difference e(t) supplied from the adder 950. Namely, the TEQ 90 and the target channel block 940 carry out the adaptive operation such that the output e(t) from the adder 950 becomes “0.”\nAs a result, the TEQ 90 has such a processing characteristic as to place ISI within only the cyclic prefix of 16 samples, as shown in FIG. 8(E).\n(2-3) Adaptive Operation Algorithm of FEQ\nNext, description will be made of an example of an adaptive operation algorithm of the above FEQ 120.\nAs described above, the FEQ 120 is an equalizer in frequency domain, which equalizes effects (changes in amplitude characteristics and phase characteristics) on a plurality of carriers having different frequencies exerted when the carriers pass through the metallic line 70 such as to equalize the characteristics of all the carriers. For this purpose, the FEQ 120 is provided with circuits shown in FIG. 10 in number corresponding to the number of the above carriers. Incidentally, identical numerals in FIG. 10 designate corresponding parts in FIG. 6.\nIn FIG. 10, a coefficient unit 1010 multiplies an input signal (received signal) Yi by a coefficient Wi. A value of the above coefficient Wi can be changed according to an output of adder 1020. In FIG. 10, a decision unit (Decision) 1030 detects a decision value Xi that is assumed to correspond to the input signal Yi, and outputs it. The adder 1020 computes a difference Ei between an output Xi of the decision unit 1030 and an output Zi of the coefficient unit 1020, and sets a coefficient Wi of the coefficient unit 1010 according to a result of the computation.\nIn the circuit structured as above, a decision value Xi is determined by inputting an output Zi of the FEQ 120 to the decoder 130, a difference Ei between the decision value Xi and an output Zi of the FEQ 120 is determined by the adder 1020, and a coefficient Wi of the coefficient unit 1010 is such adjusted that the difference Ei becomes “0.” As a result, the above equalizing process in frequency domain is adaptively performed. Incidentally, the above decoder 130 converts a decision value Xi into a bit sequence bi, and supplies the bit sequence bi to the parallel to serial buffer 140 shown in FIG. 6.\n(3) Crosstalk from ISDN Ping-Pong Transmission Line\nSince ADSL is a technique using communication lines, there are some instances where an ISDN ping-pong transmission line [TMC (Time Compression Multiplex) line] 70′ exists in the vicinity of the metallic line 70 (hereinafter referred to as an ADSL line 70) (in concrete, the ADSL line 70 and the TCM line are laid in parallel in the same cable). In such case, the ADSL line 70 is affected by crosstalk (TCM crosstalk) from the TCM line 70′.\nIn TCM transmission, in synchronization with a signal [TTR (TCM-ISDN Timing Reference)] 310 at 400 Hz as shown in, for example, FIG. 12(A), an office ISDN machine 630 shown FIG. 11 transmits downstream data in the first half cycle of the TTR 310, and a subscriber ISDN machine 640 transmits upstream data in the latter half cycle. For this, the office ADSL machine 605 is affected by near-end crosstalk (NEXT) 320 from the downstream data of the office ISDN machine in the first half cycle of the TTR 310, and affected by far-end crosstalk (FEXT) 330 from the upstream data of the subscriber ISDN machine 640 in the latter half cycle, as shown in FIG. 12(B).\nIn contrast to the office ADSL apparatus 650, the subscriber ADSL machine 660 is affected by FEXT 340 from downstream data of the office ISDN machine 603 in the first half cycle of the TTR 310, and affected by NEXT 350 from upstream data of the subscriber ISDN machine 640 in the latter half cycle, as shown in FIG. 12(C) Hereinafter, a noise section affected by NEXT will be referred to as an NEXT section, and a noise section affected by FEXT will be referred to as an FEXT section. Incidentally, FIG. 12(D) shows the NEXT section and the FEXT section in the subscriber ADSL machine 660. The effect in the NEXT section is generally larger than the effect in the FEXT section.\n(4) Description of Sliding Window\nIn order to transmit well ADSL signals in the environment where TCM crosstalk (NEXT section, FEXT section) mentioned above exists, “sliding window” is introduced. The “sliding window” is used to specify the FEXT section in which the effect of TCM crosstalk is smaller than that in the NEXT section. Effective use of the specified FEXT section can minimize the effect of TCM crosstalk, thereby certainly transmitting data.\nIn the downstream direction where an ADSL signal is transmitted from the office ADSL machine 650 to the subscriber ADSL machine 660, states of the ADSL signal are defined as follows using the “sliding window”:\nAs shown in FIG. 12(E), when a transmit DMT symbol 360 is completely involved in the FEXT section in the subscriber ADSL machine 660, the office ADSL machine 650 transmits the symbol as an FEXT symbol using a sliding window 370. When the transmit symbol is at least partly involved in the NEXT section of the subscriber ADSL machine 660, the office ADSL machine 650 transmits the symbol as an NEXT symbol. Such transmitting method is called a dual bit map method.\nIn the upstream, the subscriber ADSL machine 660 transmits the DMT symbol in the similar manner. In the downstream, the office ADSL machine 650 may transmit a pilot tone instead of the NEXT symbol in the NEXT symbol section. Such transmitting system is called an FEXT bit map system. In the FEXT bit map system, the subscriber ADSL machine 660 transmits nothing in the downstream in a section of the NEXT symbol.\n(5) Description of Frame Structure in ADSL Communication\nThe above “sliding window” is asynchronous with the TTR in TCM transmission. Here, a hyperframe used in ADSL communication in synchronization with the TTR will be explained.\nIn ADSL communication, one frame corresponds to one symbol. In steady communication, 69 frames, which are 68 ADSL frames 410 for user data and one frame for a synchronization symbol S, form one superframe 420, as shown in FIG. 13(C), for example. Further, as shown in FIG. 13(B), five superframes 420 form one hyperframe 430.\nInto the hyperframe 430, an inverse synchronization symbol I may be inserted instead of the synchronization symbol S. The inverse synchronization symbol I corresponds to a signal whose phase is rotated 180 degrees from a phase of the synchronization symbol S in the case of carriers other than the pilot tone, as shown in FIG. 14(B). In the case of the pilot tone, the inverse synchronization symbol I corresponds to a signal having the same phase as the synchronization symbol S, as shown in FIG. 14(A).\nIn the case of a hyperframe in the downstream from the office ADSL machine 650 to the subscriber ADSL machine 660, the inverse synchronization symbol I is decided to be placed (inserted) in the fourth superframe 420 in one hyperframe 430, as shown in FIG. 13(B). In the upstream, the inverse synchronization symbol I is involved in the first superframe 420 in one hyperframe 430. As shown in FIG. 13(A), one hyper frame 430 synchronizes with 34 cycles of the TTR 310 in the above-mentioned TCM transmission.\n(6) Report on TTR 310 Phase Information to the Subscriber ADSL Machine 660.\nIn order that the above mentioned ADSL machine 650 or 660 transmits and receives data using the sliding window 370 and the hyperframe 430, it is necessary for the ADSL machine 650 or 660 to recognize in which frame in the hyperframe 430 data now being transmitted/received in synchronization with the TTR 310 is, or whether the data is transmitted/received as an FEXT symbol or an NEXT symbol.\nAs shown in FIG. 11, the ADSL machine 650 on the office's side 610 can easily obtain phase information 670 on the TTR 310 from the ISDN machine 630 when installed on the same office's side 610 as the office ISDN machine 630. However, the subscriber ADSL machine 660 is separately installed in the subscriber's premise (on the subscriber's side 620), in general, so that the subscriber ADSL machine 660 cannot obtain the phase information 670 on the TTR 310 from the ISDN machine 640 installed in other subscriber's premise. For this, the subscriber ADSL machine 660 is required to receive the phase information 670 on the TTR from the office ADSL machine 650.\nThe office ADSL machine 650 therefore transmits the phase information 670 on the TTR 310 received from the office ISDN machine 630 using a carrier for reporting to the subscriber ADSL machine 660, at the time of initialization before the ADSL communication starts. Namely, the office ADSL machine 650 transmits the phase information 670 on the TTR 310 as the FEXT symbol [refer to FIG. 15(A)] or the NEXT symbol [refer to FIG. 15(B)] as decided by the sliding window 370 described above with reference to FIG. 12(E). At this time, the FEXT symbol and the NEXT symbol differ from each other in only phase.\nWhich frame in the hyper frame 430 should be transmitted as the FEXT symbol, or whether the frame should be transmitted as the NEXT symbol is decided in relation to the sliding window 370. Accordingly, the subscriber ADSL machine 660 can recognize which frame in the hyperframe is now received by receiving the above signal (FEXT symbol or NEXT symbol) transmitted from the office ADSL machine 650.\nSince one hyperframe 430 synchronizes with 34 cycles of the TTR 310, as stated above, the subscriber ADSL machine 660 can obtain the phase information 670 on the TTR 310 from the above frame position information in the received hyperframe 430, thereby synchronizing with the TTR 310. After synchronization with the TTR 310 is established, the synchronization is kept using the above pilot tone.\nAs above, when the TCM line 70′ exists in the vicinity of the ADSL line 70, the ADSL machines 650 and 660 transmit/receive data using the sliding window 370 and the hyper frame 430 to minimize the effect of TCM crosstalk, thereby realizing reliable data transmission.\nHowever, the above ADSL machine of “G.lite”, for example, is not provided with a splitter, that is, a low-pass filter (LPF). Accordingly, when the subscriber ADSL machine 660 cannot accurately receive a pilot tone for maintaining synchronization transmitted from the office ADSL machine 650 for a long time (1.8 ms or longer, for example) in steady communication due to effects of impulse noise caused by off-hook or higher harmonic noise caused by ringer generated in a telephone connected to the same line as the subscriber ADSL machine 660, or crosstalk noise caused by off-hook or ringer generated in a telephone on the adjacent line, the subscriber ADSL machine 660 cannot synchronize with the TTR 310. In such case, the ADSL communication thereafter cannot be continued. In order to re-communicate, it is necessary to initialize once more in the present condition. As this, the ADSL communication is interrupted for a long time once the synchronization with the TTR 310 becomes off."}
{"text": "1. Field of the Invention\nThe present invention relates to a charging circuit and, more particularly, to a capacitor charging circuit with a soft-start function, applicable to effectively charge photoflash capacitors of digital still cameras or other capacitive loads over a wide voltage range.\n2. Description of the Related Art\nIn photoflash systems of digital still cameras, a battery voltage source of approximately 3 volts is supplied to charge a photoflash capacitor through a transformer for raising the capacitor voltage from zero to several hundreds volts required for the excitation of a photoflash lamp. U.S. Pat. No. 6,411,064, U.S. Pat. No. 6,417,649, U.S. Pat. No. 6,518,733, and U.S. Pat. No. 6,636,021, all incorporated herein by reference, have already disclosed a variety of capacitor charging methods and circuits for charging capacitive loads such as photoflash capacitors over a wide voltage range. However, several problems are encountered in the prior art charging circuits and still need to be cured with further improvements and innovations.\nAt first, during an initial period of a charging process, a large amount of inrush current flows into a primary winding of the transformer, resulting in a significantly large amount of energy to be stored in the transformer. However, it takes an extremely long time for delivering the stored energy to the capacitive load through a secondary winding of the transformer since the terminal voltage across the capacitive load is almost zero at that time. As a result, the charging process is adversely affected by a very low efficiency during the initial period.\nEach time when the power switch coupled to the primary winding performs ON/OFF switching operations, both of the primary winding current and the secondary winding current inevitably fluctuate up and down with a great proportion to their DC mean values. Because the primary and secondary winding currents are detected in order to determine the ON/OFF timing of the power switch, such fluctuant noise that casts an incredulous shadow over the current detecting results may destroy the control mechanism of the charging circuit.\nAfter continuously supplying energy in the photoflash system for a long time, the battery voltage source may be subjected to a significant drop. The lower the battery voltage, the longer the primary winding current increases to a predetermined upper limit which controls the ON-to-OFF switching operation of the power switch. Consequently, each charging cycle is prolonged, resulting in large output ripples and low charging efficiency."}
{"text": "The present invention relates to a data transmission interface for transmitting data between a system and an utilization device, such as a printer, and more particularly, to an interface which utilizes the system's power connection to effect the data transmission.\nVarious data transmission interfaces exist in the prior art for transmitting information from a system unit, such as a computer, to an utilization device, such as a printer. These interfaces include both parallel and serial data transmission and usually consist of a separate dedicated channel over which the information is transmitted. Both the system unit and the printer unit are often powered by differing power sources connected directly thereto. Further, it is known to transmit information over a power distribution system such as that found in a household or office. Such systems are employed in conjunction with microprocessors to control various household devices, such as burglar alarms.\nPatient monitor devices are also known which operate to analyze patient samples and provide an output indication of the results of the analysis on a temporary display device. These devices are generally small allowing for portability so that they can be placed near a patient and thus do not incorporate a self-contained mechanism for creating a permanent record such as a printer. An example of such a patient monitor is the Biotrack 512 monitor. This monitor may be disconnected from its power source and operate on battery power while temporarily being located adjacent to patient in order to collect and analyze patient samples. The monitor may then be plugged back into its power source in order to recharge its batteries or perform analysis on patient samples brought to the monitor. It is important that simple connections be established between the monitor and the utilization device in order to prevent setup error as well as to conserve space and thus to make the system readily portable.\nWhile the above summarized prior art is useful in providing data communication between devices, it does not suggest an approach which relies on a simple connection of both power and signal communication facilities between a portable monitor and an utilization device, such as a printer, operating in an environment where communication must be error-free in order to insure accuracy in the communicated data."}
{"text": "A typical household utilizes numerous appliances for functioning in everyday life, and many of these appliances require electricity to operate. Most electrical appliances come equipped with a power cord that is plugged into an outlet in order to receive the necessary power. Frequently, when an electrical appliance is not in use, the power cord is wrapped up (often in a rather haphazard fashion), and the appliance is put away into a storage area. However, this ordinary process has several disadvantages.\nFirst, it is inconvenient to repeatedly remove an item from storage for use, to plug it in to an electrical outlet, and then to unplug and return it again. This is not to mention the difficulty of wrapping up and storing power cords, which are built rather inflexibly and often end up in only a tangled mess. Leaving all appliances out of storage and always available for use might be an option, but besides the lack of adequate space, there are often not enough outlets to accommodate powering multiple items simultaneously. Even when there are enough outlets, many power cords coming from multiple appliances can be a safety hazard and an eye-sore.\nTherefore, it is desirable to have a quick and easy way to connect and disconnect an appliance to and from power without the use of a power cord. Furthermore, convenient storage and adequate space problems could be solved by a system that allowed a single mechanism to simultaneously supply electrical power and also support an electrical appliance, such as by suspending the appliance from a wall.\nSome designs, such as powered rails, allow for simultaneous transfer of electrical power and structural support to appliances. However, these designs often leave the electrical conductors exposed, so that safety concerns may offset the advantages of speed and simplicity achieved through the lack of a power cord. Attempts to improve the safety of powered rails through the addition of insulative materials have complicated the connecting process. Using insulation often requires precise alignment of the electrical conductors of the powered rails with electrical conductors of the appliances, requiring more time and effort. Additionally, forming a connection between the two sets of conductors often requires greater force in order to move the electrical conductors of the appliance through the insulation to the electrical conductors of the powered rails. A similarly greater force is then required to disconnect the appliance from the powered rails. In such a case, a tradeoff is made sacrificing ease and simplicity for safety. Other connectors that provide both electrical power and structural support for appliances make similar tradeoffs.\nTherefore, in light of the foregoing, what is needed is a mechanism that provides a quick and easy means for providing both electrical power and structural support to an appliance, but a means that is safe and that requires minimal physical force."}
{"text": "The present invention relates generally to a lens for sunglasses, and more particularly to a uniquely configured quadroradial lens having a dual radius horizontal curvature and a constant radius vertical curvature for minimizing the prismatic shift of the lens. As discussed in greater detail below, embodiments of the present invention maximize the interception of peripheral light, while, at the same time, mitigating prismatic shift to enhance optical resolution for the wearer. An optional frusto-conical as worn configuration of the lens also permits construction of sunglasses which conform closely to the front and sides of the wearer's head. The resulting low profile glasses utilizing the lens of the present invention are particularly suited for demanding situations which require precise optical resolution and interception of peripheral light, such as competition skiing or bicycle racing.\nSunglasses have long been designed with the general objective of blocking the sun or other sources of bright light from one's eyes. Initially, numerous designs of dual lens glasses were developed, differing essentially only in aesthetic features. However, the unitary lens was later developed, and together with existing dual lens designs, has been geometrically modified in response to various optical considerations such as optical clarity and resolution, field of vision, light wave refraction, and others.\nAlthough prior dual lens designs are useful for some purposes, the conventional dual lens system is inherently incapable of meeting the demands of certain activities. For example, the frame on dual lens glasses presents a substantial obstruction to one's peripheral vision, which can be extremely disadvantageous in any fast-paced activity. Even in dual lens glasses without a frame along the lower edges of the lenses, the edge of each lens disrupts peripheral vision. Simply providing a larger lens that is typically stamped or molded from a flat plane or a spherical blank, causes the glasses to extend too far tangentially away from the side of the head, leaving the glasses with an undesirably bulky profile.\nAt the same time, conventional dual lens glasses only intercept sunlight directly in front of the eye, leaving a large, unprotected periphery about each lens. Momentary flashes of light around the lens during activity cause constriction of the pupils, with a fleeting blindness as one attempts to readjust through the darkened lens.\nPrior art attempts to block this peripheral light (aside from opaque blinders) included bending (e.g., thermoforming) a flexible lens in a posterior direction near the lateral edge. Although this improved the interception of peripheral light, the resulting optical resolution was unacceptable for high speed competition situations. This is due to the phenomena that even minor irregularities in the radius of curvature, which inherently result when bending a lens, cause an irregular diffraction of light waves passing through that region of the lens and distort the field of vision.\nA unitary, molded, frusto-conical lens blank was then developed, such as that disclosed in U.S. Pat. No. 4,515,448 issued to Tackles. The frusto-conical lens was designed to conform closely to the wearer's head, having lateral edges that curved and extend interiorly to block peripheral light that would otherwise pass around the lenses of a dual lens design. In addition, diffraction gradients common in the bent, flexible lens, were minimized by molding the frusto-conical lens with a predetermined curvature. Nevertheless, the potential for improvement remained for several reasons.\nThe next improvement in specialty unitary eyewear utilized a unitary cylindrical lens which was curved about an axis having a substantially constant radius throughout, such that the lens defined a portion of the wall of a cylinder. This lens is the subject of U.S. Pat. No. 4,859,048, issued to Jannard. The cylindrical lens demonstrated improved optical properties and interception of peripheral light. Further, it provided a sleek, low profile design that also improved ventilation. Nevertheless, the cylindrical lens produced a measurable prismatic shift in each of the vertical and horizontal planes, particularly at off axis viewing angles.\nAnother improvement in specialty unitary eyewear utilized a toroidal lens that was curved along each of two substantially perpendicular axes to produce a lens of generally toroidal configuration. This lens is the subject of U.S. Pat. No. 4,867,550, also issued to Jannard. The toroidal lens configuration was characterized by having arcuately configured horizontal and vertical cross-sections. The arc of each respective cross-section corresponded to a given radius. Thus, the toroidal lens can be defined by two different radii, and is configured to substantially conform to the wearer's head. In addition, the lens thickness was tapered to further reduce prismatic shift.\nYet another improvement in unitary lens eyewear utilized an elliptical lens. This lens is the subject of U.S. Pat. No. 5,774,201, also issued to Tackles. The elliptical lens is characterized as having a horizontal cross-section that substantially conforms to an elliptical shape. The vertical cross-section of the elliptical lens can be defined by any of various geometric shapes, such as a cylinder, a cone, an ellipsoid, or an ellipsoid of revolution. The elliptical lens was introduced to improve the optical characteristics of specialty eye wear, including refraction between medial light entering at the front of the lens and peripheral light entering at the lateral ends of the lens.\nNotwithstanding the many advantages presented by this progression of unitary lens designs, there is a continuing need for a specialty lens having excellent optical qualities and providing reduced light wave defraction of medial, lateral and peripheral light, while at the same time providing a low profile, aerodynamic configuration that allowed for adequate ventilation and maximum comfort. Preferably, the lens will exhibit good impact resistance, as well as minimal prismatic shift in both the horizontal and vertical axes."}
{"text": "Diesel engine builders have been developing new engines with advanced emissions controls to meet greatly tightened EPA emission requirements for engines due to go into effect in the year 2007. Effective with the 2007 model year, the U.S. Environmental Protection Agency set stringent limits on nitrogen oxide (NOx) and particulate matter (PM) emissions from on-highway trucks and buses. Engines incorporating diesel oxidation catalysts and particulate filters will require significantly revised lubricant formulations than in the past.\nConsequently, engine manufacturers have been developing diesel engines that not only use diesel particulate filters that trap and further reduce soot emissions but also which operate on pollution-reducing ultra-low-sulfur diesel(ULSD) fuel and utilize cooled exhaust gas recirculation (EGR) devices to redirect some of the exhaust gases normally emitted by the vehicle back into the engine, creating more internal soot.\nBecause the new engines operate at lower combustion and higher coolant temperatures and do not burn off the soot and other particles, and because the after treatment devices are easily damaged, existing oils, even those meeting the CI-4 or CI-4 PLUS standard, do not provide the necessary protection to meet the new emission requirements for on-highway trucks and buses.\nIn particular, certain performance enhancing additives used in the past will cause problems in their currently used concentrations in the new generation of diesel engines which utilize relatively sophisticated exhaust filtration systems. For instance, ash and other metallic components of the traditional additives are incombustible, and thus would remain in the system and threaten to foul the sophisticated filters that will be on the new engines, increasing maintenance costs and adding to the emissions. Therefore, low sulfur-ash-phosphorus (“low SAP”) lubricant oils will be required in the industry which incorporate limits placed on certain additive components in the lubricant oil as follows: <1.0% ash, <0.12% phosphorous, and <0.4% sulfur.\nThe American Petroleum Institute (API) has promulgated a new lubricant category designated “CJ-4,” that is specifically engineered to aid the next generation of low-emission diesel engines in meeting the 2007 engine emission standards. In order to comply with the new API CJ-4 standard, lubricant formulators and manufacturers have been facing the challenge of developing new cleaner lubricant additives and components that can deliver the combination of performance and purity that the new engines and standards demand.\nFor off-road use diesel engine applications, other pre-existing API oil specifications will still apply as applicable. For example, oils in the API CI-4 category are formulated for use in high-speed, four-stroke cycle diesel engines and were designed to meet the EPA exhaust emission standards that came into force in October 2002. These oils have been compounded for use in all applications with diesel fuels ranging in sulfur content up to 0.05% by weight. CI-4 PLUS, adopted in 2004, represented an enhancement to the CI-4 category. Oils meeting the CI-4 PLUS designation are formulated to meet additional soot-control requirements of specific engines with cooled EGR.\nIn addition to meeting ever-tightening industry standards for oil performance and purity, the present investigators also have identified a coolant system leak problem that can occur in diesel engines, such as small leaks in seals of engine coolant systems used in EGR engines, which can threaten the performance and proper operation of lubrication system. These small coolant leaks can lead to contamination of the engine lubricating oil, viz., the crankcase oil, with coolant (i.e., typically ethylene glycol, propylene glycol and/or other (poly)alkylene glycol and the like). Coolants may also contain corrosion inhibitors, and other additive agents (including bittering agents, etc.) The presence of the coolant in the lubricant system has been observed by the present investigators to induce and/or promote plugging problems in the oil filter.\nMany patents have taught the manufacture and use of dispersant VII materials in engine oil formulations. These patents include, for example, U.S. Pat. Nos. 6,107,257; 5,182,041; 5,188,745; 5,238,588, 4,8636,623; 5,075,383; 6,107,258; 5,556,923; and European patents EP 0922752; EP 0909805; EP 0491456; EP 0510892; EP 0338672EP 0396297; EP0549196; and EP 0417904, all of which are incorporated herein by reference in their entirety for their teaching. None of these patents are understood to recognize and address both CJ-4, CI-4 Plus and CI-4 compliance and coolant contamination issues for diesel engine lubricating compositions.\nTherefore, the present investigators have identified a need for CJ-4, CI-4 Plus and CI-4 lubricant formulations that can further multi-task to minimize or prevent oil filter plugging in heavy-duty diesel engines, such as cooled EGR engines. As will become apparent from the descriptions that follow, the invention addresses this need as well as providing other advantages and benefits."}
{"text": "This invention relates to a unique motor generator set. In particular to a motor generator set where in the output of a pair of generators is fed back to a battery as well as directly to an electric motor.\nIn the past motor generator sets have been used in automobiles where the electric generated by the generator is sent back to the battery. This only resulted in charging the battery."}
{"text": "Flat panel display devices are increasingly gaining market acceptance for a variety of different applications. For example, active matrix liquid crystal displays (AMLCD's) have found widespread use as the video monitors in laptop computers, video cameras and avionic navigation modules, to name but a few devices. Other types of display devices such as electroluminescent (EL) and field emissive displays (FED's) are also used in a variety of industrial and consumer applications. The advantage of each of these types of devices resides in the fact that they are all substantially flat, particularly as compared to the cathode ray tube that has been in use for the past fifty years.\nIn the AMLCD, the elements which cause the device to effect a desired optical characteristic are typically sandwiched between a pair of thin glass plates. These elements include first and second patterned electrodes for applying an electrical field to liquid crystal (LC) material disposed therebetween. Each pair of oppositely disposed patterned electrodes define a single picture element or pixel. The liquid crystal material typically is a conventionally known liquid crystal material, such as cholesteric, polymer dispersed liquid crystal materials, twisted nematic (TN), supertwist nematic (STN), chiral smectic and others. The application of an electrical field to the LC material causes it to change its orientation from a first condition to a second condition, for example, transparent to opaque. However, in order to control the orientation of the liquid crystals, it is necessary to provide numerous other optical elements, such as at least a pair of polarizers, and a plurality of alignment layers. A conventional AMLCD is fully described in, for example, U.S. Pat. Nos. 4,666,252, 4,715,685 and 5,061,040 all to Yaniv, et al., the disclosures of which are incorporated herein by reference. An additional U.S. Pat. No. 4,961,630 to Yaniv and Baron, teach an AMLCD having three electrodes, the third electrode provided to increase device capacitance.\nUnfortunately, the construction of conventional AMLCD's and STN based displays generally is that using a twisted mode configuration leads to numerous deficiencies and disadvantages. For example, the need to provide two polarizers for conditioning the optical output substantially lowers the transparency of the device. The result is a darker display or alternatively one requiring a larger, i.e., higher powered backlight Accordingly, the polarizers either result in the need for larger backlights, increasing cost or have poorer color intensity, resulting in diminished display performance. Additionally, at least one, and typically a pair of alignment layers are necessary for purposes of properly orienting the liquid crystal molecules upon the application or removal of the electrical field. However, these layers must be carefully applied in order to achieve perfect orientation. The steps involved in depositing and preparing the alignment layers are difficult, time consuming, and introduce numerous opportunities for defects in the devices. Accordingly, the alignment layer contributes to lower device manufacturing yields and increased device costs.\nAdditional limitations to conventional AMLCD's resides in the basic characteristics of the LC material. Specifically, upon application of an electric field, the LC molecules will align themselves according to the field, including any effect of the interposed alignment layers, providing a desired optical effect Removal of the electrical field allows the LC molecules to \"relax\" back to the original state of orientation. However, the speed of relaxation is considerably slower than the speed of orientation in response to the electric field. This phenomenon has severe consequences for high speed operation of AMLCD's, and in particular STN's.\nThe problems noted above are further exacerbated when the device reaches higher temperatures, as can happen upon prolonged exposure to high intensity backlights. The high temperature dependency also has substantial consequences in terms of the types of application in which such a display may be used. For example, poor high temperature performance eliminates reliable use of AMLCD's in automotive applications. Operation speed is likewise deleteriously effected by lower temperatures, which substantially slow both excitation and relaxation speeds.\nAn additional deficiency of conventional AMLCD's relates to the relatively poor viewing angles of the devices. By this it is meant that the display appearance, at angles substantially off 90 degrees to the surface of the display, is substantially degraded. This degradation is due to the inherent characteristics of the polarized light emitted from a twisted configuration (IN, STN, etc.) AMLCD in conjunction with the need to interpose polarizers and alignment layers on the device substrates.\nAccordingly, there exists a need for a display device which provides the desired changes and control in optical characteristics, while avoiding the problems inherent in conventional LCD's. Such a device should be easier to fabricate, have fewer optical components, and superior optical and electrical efficiency. The devices should also be easily adaptable to conventional semiconductor fabrication techniques, or even better, to screen printing technology, which is both simpler and less costly. The improved device should reduce dependency on polarizers and alignment layers to effect the desired optical performance."}
{"text": "Eledoisin, which is isolated from an acetone extract of the posterior gland of the ocean Eledone moschata and E. aldrovandi and having a strong antihypertensive activity against a dog, is a physiologically active neuropeptide consisting of 11 amino acid residues (Erspamer, V., et al., Experientia, 18, 58, 1962). It has been obvious that the eledoisin shows a constricting activity of the ileum of a guinea pig other than an antihypertensive activity and also shows a specific action such as acceleration of salivary secretion of dog administered by intravenous injection.\nThereafter, the peptide having the same activities has been isolated from a skin of frog and this has been named physalaemin (Erspamer, V., et al. Experientia, 20, 489, 1964).\nSuch peptides were named tachykinin (tachy(=fast)kinin, something to contract quickly) against bradykinin (brady (=gradually)kinin, something to move gradually) because they have constricting activity against the ileum of a guinea pig quickly. Further, the chemical structure of substance P has been clarified based on the chemical structure of such peptides.\nThough peptides classified into the tachykinin have been isolated thereafter one after another from the Amphibian class (Yasuhara, T., et al., Biomed. Res., 2, 613, 1981) and the Avian class (Conlon, J. M., et al., Regulatory Peptides, 20, 171, 1988), sialokinin only has been found out from a mosquito which has mediated yellow fever other than eledoisin from a invertebrate animal (Champagne, D. E. & Ribeiro. J. M.: Proc. Natl. Acad. Sci., USA, 91, 138, 1994).\nAt present, tachykinin has a common amino acid sequence represented by the following formula in the C-terminal of a peptide:\nPhe-Xaa-Gly-Leu-Met-NH2 (SEQ ID NO: 6)\nwherein, Xaa is aromatic amino acid (e.g., Phe, Tyr) and branched amino acid (e.g., Val, Ile). Tachykinin is used herein as a generic term of a physiologically active peptide indicating an intestine constricting activity, a hypotensive activity, a saliva secretion accelerating activity and the like. In such kinds of tachykinin, there are included aforementioned substance P, eledoisin, physalamine, neurokinin A, neurokinin B, kassinin and the like (SEIBUTUGAKU-JITEN (Ver. 4). Iwanami Shoten, 1997).\nIn this way, it has been expected to employ the tachykinin peptide as a basic chemical compound for developing a new medicine. In particular, since it has been considered that substance P has taken part in transmitting a pain as a transmitting substance of a primary sensory nerve, it has been expected and researched to develop as analgesics. Further, it has been expected to use practically a chemical reagent for elucidating an information processing mechanism in a nerve system of higher animals.\nBy the way, it is said that the posterior salivary gland of octopuses are venom gland since some of the salivary gland may include toxic substances such as tetrodotoxin known as fugu poison, and cephalotoxin that is a glycoprotein having an effect of numbing Crustaceans other than aforementioned eledoisin. Further, it has been known that they might include a biogenic amine group such as octopamine, serotonin, tyramine, noradrenaline, histamine, and acetylcholine and they might include an enzymatic group such as a proteolytic enzyme and hyaluronidase (Boucaud-Camou, E. & Boucher-Rodoni, R. 1983. Feeding and digestion in Cephalopods. In “The mollusca”, (Saleuddin, A. S. M. & Wilbur, K. M.), Academic press and New York). It, however, has not been reported that the tachykinin peptide other than eledoisin has been found. Further, though it has been discovered that eledoisin shows smooth muscle constricting activity, angiectatic activity, and antihypertensive activity to the mammals, it has not been discovered that eledoisin has a role to the octopus.\nIn this way, it is required to find out the tachykinin peptides from greater number of animal species, in order to obtain information which is useful for development of the medicines and for solving an information processing mechanism in a nerve system of higher animals based on research and the like of a structure-activity relationship by elucidating a species specificity of the peptides and a role of each of animal species about the tachykinin peptides.\nTherefore, the object of the present invention is to provide the novel tachykinin peptide capable of being used as a chemical reagent for elucidating an information processing mechanism in a nerve system of higher animals and as a basic chemical compound in developing the medicines and the pesticides by finding out the novel tachykinin peptides, making clear the structure thereof, and further elucidating the physiologically activities thereof.\nFurther, the object of the present invention is to provide a method for producing the tachykinin peptide by identifying the precursor polypeptide of such tachykinin peptide and gene encoding thereof."}
{"text": "Automatic step systems for recreational vehicles, motor homes, and the like are well known in the art. These systems are typically electrically-controlled and electrically-actuated to extend and retract an entryway step in response to a signal provided by an individual wishing to enter or exit the vehicle. One common system extends the step when the vehicle door is opened, and then retracts the step when the vehicle door is closed. Other systems offer a switch located just inside the vehicle door that controls extension and retraction of the step.\nOther systems incorporate a motor assembly for automatically extending and retracting the step assembly. The motor rotates a pivot rod through a gear assembly that is coupled to the rod. The pivot rod moves a linkage assembly to extend and retract the steps. However, an applied load tends to move the steps towards the retracted position. Generally, the motor assembly is used to resist the tendency of the steps to move towards the retracted position. This may limit the service life of components in the motor assembly. Other types of systems use solid stops that contact one or more of the links to resist applied loads instead of the motor assembly. However, the links do not positively lock against the solid stops. Therefore, an improved mechanism for extending and retracting steps in recreational vehicles is needed."}
{"text": "1. Field of the Invention\nThe present invention relates to reflective warning devices which are useful for alerting motorists to the presence of other vehicles and road hazards.\n2. Description of the Prior Art\nThe prior art has previously provided warning devices of both an active and passive nature which act to warn of the presence of a dangerous condition, particularly a road hazard associated with moving or stationary vehicles. Active visual devices particularly include lights disposed on a vehicle or on a stationary object intended to warn of a dangerous condition, it having been previously appreciated that a flashing light is more effective in alerting observers to the presence of potential dangers. However, active visual devices must be operated by a source of power such as a battery or the like which causes the light to repetitively flash on and off. As is obvious, such active devices require a substantial amount of energy for operation and are also costly. Passive visual warning devices, particularly reflective devices wherein light from vehicle headlights or other sources are reflected to observers, are also known in the art. Such reflective devices have even been caused to provide a flashing effect by intermittently obscuring at least portions of a reflective surface. For example, LaLonde, in U.S. Pat. No. 3,528,721, causes a flashing effect of a reflective device by mounting the device against the wheel of a bicycle so that rotation of the wheel causes the reflective device to oscillate. Priest, in U.S. Pat. No. 3,551,024, discloses a flashing retroreflective reflector in which a convex lens focuses incoming light on a retroreflective medium, the device being oscillated such that retroreflection is obtained only when incoming light is incident along the axis of the lens. Hammer, in U.S. Pat. No. 2,869,424, provides a passive warning signal device wherein portions of a reflective surface are intermittently obscured by spring mounted masking plates, the masking plates being oscillative due to motion of the vehicle carrying the reflective device. Klaenhammer, et al., in U.S. Pat. No. 4,023,888, provides a reflective structure having portions thereof intermittently obscured by a swinging masking member, the swinging member being driven by a motor. However, the prior art has not provided a passive visual warning reflector device capable of producing a \"flashing\" visual signal on mounting of the device to either a stationary or potentially movable object. A free-swinging portion of the present device has means formed thereon which are actuable by air currents to be operable both when the structure on which the device is mounted is moving and also when the structure is stationary, as long as air is circulating about the device when the structure is stationary."}
{"text": "In general, semiconductor device manufacturers conduct testing for device reliability in the evaluation of new products, design changes, and process changes. Most manufacturers subject their parts to a short period of accelerated stress testing prior to shipment in an effort to eliminate infant mortality failures. There is also a serious effort with many manufacturers to eliminate the infant mortality phenomena by some form of stress testing during wafer sorting. Even after the introduction of a new product or package, on-going monitoring of reliability and quality normally continues throughout the production life cycle. On-going reliability testing is normally performed both in order to monitor device performance and to accumulate statistical information.\nAs a result of such testing, possible wear-out mechanisms are identified and eliminated, either through modifications in the process or in the design. However, the manufactured devices may continue to exhibit early life failures or infant mortality failures. Examples of such defects includes oxide pinholes, photoresist or etching defects that cause near-opens or shorts, and contamination on the chip or in the package. A number of stress tests have been developed to accelerate the effects of various failure mechanisms. By subjecting a device to extreme operating conditions for a short period of time, such as for example, putting a device into a condition of abnormally high temperature and high voltage (such a test is called \"burn-in test\"), one can with some specified confidence level predict life and device performance under more normal conditions. In such a burn-in test, a static super voltage (about 7.0 to 8.5 volts) is applied to one or more pins of a device. Integrated semiconductor memory devices such as DRAMs and SRAMs typically include various input buffers for receiving externally provided address signals, data signals, and control signals via their pins (or pads). Each buffer acts as an input stage to circuits within the memory devices.\nFIGS. 1 and 2 illustrate examples of the most widely used MOS input buffers for semiconductor devices. Referring first to FIG. 1, the CMOS input buffer 50 includes two serially-connected inverters 6 and 12. The inverter 6 is composed of a p-channel MOS (hereinafter abbreviated as PMOS) transistor 10 and an n-channel MOS (abbreviated as NMOS) transistor 11. The complementary MOS (abbreviated as CMOS) buffer 50 receives an externally applied input signal INE and converts the external signal into a compatible internal signal INI (or INI). The gates of the CMOS transistor pair 10 and 11 are supplied with the input signal INE. At the drain junction of the CMOS transistor pair 10 and 11, the internal signal INI compatible with the external signal INE is produced.\nReferring now to FIG. 2, the CMOS input buffer circuit 60 includes a NOR gate 8 and an inverter 17. The NOR gate 8 is formed of PMOS transistors 13 and 14, and NMOS transistors 15 and 16. An input signal INE is externally applied to the gates of the transistors 14 and 15. An external control signal such as chip select signal CS is applied to the gates of the transistors 13 and 16.\nGate oxide is an important element of MOS transistors. This thin dielectric layer can break down, resulting in gate shorts, during a long or very strong application of electric field across the oxide. Oxide breakdown is generally believed to be caused by positive charge buildup. Therefore, in a reliability test mode such as burn-in test, input buffer oxide breakdown of the targeted device may occur due to the high voltage which is applied to the device for reliability evaluation. More specifically, in the input buffer 50 of FIG. 1, the NMOS transistor 11 is turned on when a super voltage input signal is applied to the input terminal of the CMOS inverter 6, but the PMOS transistor 10 is turned off. The drain voltage of the PMOS transistor 10 is then about Vss (i.e., 0 volt) while the source voltage thereof is Vdd. The difference of gate-source and gate-drain voltages of the PMOS transistor 10 may exert stress on its gate oxide and cause gate oxide breakdown. Similarly, in such an input buffer 60 of FIG. 2, there is a high probability that the gate oxide of the PMOS transistor 14 will be broken down by the stress due to the large difference between its gate-drain and gate-source voltages."}
{"text": "Orally administered medicaments are given to the patient in many forms, such as liquid solutions, emulsions, or suspensions, or in solid form such as capsules or tablets (as used herein, the term \"tablet\" means any shaped and compressed solid dosage form, including caplets). Medicaments administered in tablet or capsule form are usually intended to be swallowed whole. Therefore, the often disagreeable taste of the active ingredient need not be taken into account in formulating the medicine, except for the provision of means to prevent the taste from being apparent during the short time that the medicine is in the mouth. Such means may include the provision of an appropriately thin and quickly dissolving coating on the tablet, the use of the gelatin capsule form (the gelatin outer shell of the capsule keeps the active ingredient inside until the capsule has been swallowed), or simply compressing a tablet firmly so that it will not begin to disintegrate during the short time that it is intended to be in the mouth.\nChildren, older persons, and many other persons have trouble swallowing whole tablets and even capsules. Therefore, in cases where the dosage to be administered cannot be made into a very small tablet or capsule, it is desirable to provide the medicine either in liquid form or in a chewable solid form, in addition to the tablet or capsule that is designed to be swallowed whole. Even where the medicine can be formulated as a liquid, it is desirable also to be able to provide a chewable solid form (i.e. tablets) because of added convenience versus carrying a supply of liquid medicine.\nA common problem with chewable tablet forms is the often disagreeable taste of the active ingredient which manifests itself during chewing. In some cases, the taste of the active medicament in a tablet can be overpowered by adding flavoring ingredients to the tablet so that when it is chewed, the taste of the active ingredient is simply overpowered. For instance, this has been done with children's aspirin where the dosage is small enough so that the amount of flavoring agents needed to mask the taste of the medicine is not so great that the tablet becomes unreasonably large. A different approach is taken with a commercially available children's size tablet of acetaminophen (acetyl para-aminophenol or \"APAP\") wherein the APAP is present in granules that are coated with ethyl cellulose. A significant proportion of the APAP remains shielded by the coating (and therefore does not contribute to taste) while the tablet is in the mouth, despite some breakage of the ethyl cellulose coating during compression of the tablet and some additional breakage of the coating during chewing. The APAP becomes bioavailable via permeation through the coating (although ethyl cellulose is not soluble in aqueous fluids, water does permeate through the coating) and from the granules where the coating is broken.\nExamples of taste masked coating systems are disclosed in the following references. U.S. Pat. No. 4,851,226, issued Jul. 25, 1989, discloses chewable medicament tablets wherein granules of active ingredient are directly coated with a blend of cellulose acetate or cellulose acetate butyrate and polyvinylpyrrolidone. Commonly assigned U.S. Pat. No. 5,215,755, issued Jun. 1, 1993, discloses chewable medicament compositions comprising a rotogranulation blend of from about 88 to about 97.5% medicament, about 2 to about 10% polyvinylpyrrolidone (PVP) and about 0.5 to about 2.0% sodium lauryl sulfate (SLS), by weight of the weight of the total composition. In further embodiments a coating of hydroxyethyl cellulose (HEC) or a mixture of hydroxyethyl cellulose and hydroxypropyl methylcellulose (HPMC) is added to these rotogranulated particles. Commonly assigned U.S. Pat. No. 5,075,114, issued Dec. 24, 1991, discloses a chewable medicament comprising a coating for active medicament comprising a polymer blend of cellulose acetate and/or cellulose butyrate and water soluble hydroxypropyl cellulose to provide a taste masked and/or sustained release coating.\nThe present invention is directed to the discovery of a reverse enteric coating process for active medicaments which can achieve a better balance between taste masking, dissolution and rate of bioavailability than other previously known coating combinations."}
{"text": "A variety of surgical instruments include a tissue cutting element and one or more elements that transmit radio frequency (RF) energy to tissue (e.g., to coagulate or seal the tissue). An example of such an electrosurgical instrument is the ENSEAL® Tissue Sealing Device by Ethicon Endo-Surgery, Inc., of Cincinnati, Ohio. Further examples of such devices and related concepts are disclosed in U.S. Pat. No. 6,500,176 entitled “Electrosurgical Systems and Techniques for Sealing Tissue,” issued Dec. 31, 2002, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,112,201 entitled “Electrosurgical Instrument and Method of Use,” issued Sep. 26, 2006, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,125,409, entitled “Electrosurgical Working End for Controlled Energy Delivery,” issued Oct. 24, 2006, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,169,146 entitled “Electrosurgical Probe and Method of Use,” issued Jan. 30, 2007, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,186,253, entitled “Electrosurgical Jaw Structure for Controlled Energy Delivery,” issued Mar. 6, 2007, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,189,233, entitled “Electrosurgical Instrument,” issued Mar. 13, 2007, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,220,951, entitled “Surgical Sealing Surfaces and Methods of Use,” issued May 22, 2007, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,309,849, entitled “Polymer Compositions Exhibiting a PTC Property and Methods of Fabrication,” issued Dec. 18, 2007, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,311,709, entitled “Electrosurgical Instrument and Method of Use,” issued Dec. 25, 2007, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,354,440, entitled “Electrosurgical Instrument and Method of Use,” issued Apr. 8, 2008, the disclosure of which is incorporated by reference herein; U.S. Pat. No. 7,381,209, entitled “Electrosurgical Instrument,” issued Jun. 3, 2008, the disclosure of which is incorporated by reference herein.\nAdditional examples of electrosurgical cutting instruments and related concepts are disclosed in U.S. Pub. No. 2011/0087218, entitled “Surgical Instrument Comprising First and Second Drive Systems Actuatable by a Common Trigger Mechanism,” published Apr. 14, 2011(now U.S. Pat. No. 8,939,974, issued Jan. 27, 2015), the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2012/0083783 (now U.S. Pat. No. 8,888,809), entitled “Surgical Instrument with Jaw Member,” published Apr. 5, 2012, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2012/0116379, entitled “Motor Driven Electrosurgical Device with Mechanical and Electrical Feedback,” published May 10, 2012 (now U.S. Pat. No. 9,161,803, issued Oct. 20, 2015), the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2012/0078243, entitled “Control Features for Articulating Surgical Device,” published Mar. 29, 2012, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2012/0078247 (pending), entitled “Articulation Joint Features for Articulating Surgical Device,” published Mar. 29, 2012, the disclosure of which is incorporated by reference herein; U.S. Pub. No. 2013/0030428, entitled “Surgical Instrument with Multi-Phase Trigger Bias,” published Jan. 31, 2013 (now U.S. Pat. No. 9,089,327, issued Jul. 28, 2015), the disclosure of which is incorporated by reference herein; and U.S. Pub. No. 2013/0023868 (pending), entitled “Surgical Instrument with Contained Dual Helix Actuator Assembly,” published Jan. 31, 2013, the disclosure of which is incorporated by reference herein.\nWhile a variety of surgical instruments have been made and used, it is believed that no one prior to the inventors has made or used the invention described in the appended claims.\nThe drawings are not intended to be limiting in any way, and it is contemplated that various embodiments of the technology may be carried out in a variety of other ways, including those not necessarily depicted in the drawings. The accompanying drawings incorporated in and forming a part of the specification illustrate several aspects of the present technology, and together with the description serve to explain the principles of the technology; it being understood, however, that this technology is not limited to the precise arrangements shown."}
{"text": "Field\nThe technology described herein relates generally to a liquid crystal display (LCD), and more particularly, to an LCD including a plurality of pixels displaying a color image.\nDiscussion of the Background\nA typical liquid crystal display (LCD) includes a pixel electrode, a common electrode, a liquid crystal (LC) layer including liquid crystal molecules configured to be tilted by an electric field that is generated by the pixel and common electrodes, and a backlight unit configured to emit light to the LC layer.\nSuch an LCD device displays a desired image by controlling intensity of the electric field generated by the pixel and common electrodes and then adjusting the transmittance of light that is emitted from the backlight unit to be transmitted through the LC layer.\nThe LC layer included in the LCD may have a phase delay characteristic that is intrinsic to the liquid crystals themselves, and accordingly, in the LCD, maximum transmittance of light transmitted through the LC layer has a characteristic according to each wavelength.\nThe above information disclosed in this Background section is only to enhancement of understanding of the background of the inventive concept, and, therefore, it may contain information that does not form the prior art that is already known to a person of ordinary skill in the art."}
{"text": "This disclosure relates to casting metal alloy materials to achieve a more uniform microstructure.\nInvestment casting is known and used for fabricating near net shape components of relatively complex geometries from high temperature alloys. Typically, the process includes forming a wax pattern of the component and coating the wax pattern with a ceramic slurry. The slurry is fired to form a refractory shell investment mold and the wax is removed from the interior to form a molding cavity within the shell. A molten alloy is then poured into the cavity to form the component."}
{"text": "1. Field of the Invention\nThis invention relates to a stepping motor having two sets of stator yoke units, and in particular to the adjustment of cogging torque.\n2. Related Background Art\nReferring to FIG. 10 of the accompanying drawings which shows the structure of stator yoke units used in a stepping motor according to the prior art, stators 1 and 2 having stator poles 1a and 2a, respectively, constructed in the form of comb teeth form a stator yoke unit. Stator yoke units constructed of entirely the same stators as the pair of stators and 2 are disposed back to back on their respective stators 2 and constitute the stator assembly of the stepping motor. The two sets of stator yoke units are such that circular bosses 2b and holes 2c are formed in each of the respective stators 2 so that the phases of the stator poles deviate from each other by 1/4 pitch.\nReferring now to FIG. 11 of the accompanying drawings which is a perspective view of the stepping motor constructed in the manner described above, the reference numeral 12 designates a rotor shaft, and a cylindrical rotor magnet 13 is fixed coaxially with the rotor shaft 12 for rotation about the rotor shaft. The rotor magnet 13 has formed on the peripheral surface thereof opposed to the stator poles 1a, 2a, 1a, 2a the same number of magnetic poles as the number of the stator poles on one set of stator yoke units. That is, the number of one set of stator poles in FIG. 11 is 24, and in this case, the number of the magnetic poles of the rotor is 24 because each of N poles and S poles is counted as one, and is 12 if counted in terms of the number of pairs of magnetic poles. The reference numerals 14 and 15 denote coils wound on the respective stator yoke units in the form of bifilar turns and led out of the stators through cut-away portions 1f, 1f, 2f and 2f.\nFIG. 12 of the accompanying drawings is a diagram of the driving circuit of the stepping motor, and FIG. 13 of the accompanying drawings is a time chart of 2-phase excitation driving pulses applied to the driving circuit of FIG. 12.\nFIG. 14a of the accompanying drawings is a developed view of stator poles and a rotor magnet of the stepping motor of FIG. 11. Two sets of stator yoke units are designated by a (comprising stators 1a.sub.1 and 2a.sub.1) and b (comprising stators 1a.sub.2 and 2a.sub.2), respectively, and the magnetic poles of the rotor are denoted by 13a.\nHere, description will be made of the stationary position of the rotor and the cogging torque when the coils are not excited.\nDescription will first be made of the cogging torque acting between the stator yoke unit a and the rotor.\nIn considering the torque acting between the stator yoke unit a comprising one set of stators and the rotor, a two-pole stepping motor shown in FIG. 15 of the accompanying drawings is adopted as a model. This two-pole stepping motor corresponds to one pitch of the stepping motor of FIG. 11.\nIf in FIG. 15, the coordinates x, y are defined as shown, the point at which the torque acting on the rotor 16 is zero is the point at which the center of the magnetic pole coincides with the center of the stator pole and therefore, .theta.=0, .pi., 2.pi., . . . , n.pi..\nAlso, when ##EQU1## the torque acts equally on the rotor 16 between the right and left stators 17 and 18 and therefore, again in this case, the torque is zero. Assuming that the torque varies in the fashion of sine, EQU Ta=a.sub.1 sin 2.theta. (1)\nand the torque varies at a period of 1/2 of the period 2.pi. of the stator poles.\nIf the number of the pairs of magnetic poles of the stator and the rotor is p, the equation (1) becomes EQU Ta=a.sub.1 sin 2p.theta.. (2)\nIf p.theta. =.alpha., EQU Ta=a.sub.1 sin 2.alpha.. (3)\nThe equation (2) is a case represented by mechanical angle, and the equation (3) is a case represented by electrical angle. Description will hereinafter be made by the use of .alpha..\nA curve Ta in FIG. 16 of the accompanying drawings represents the cogging torque acting between the stator yoke unit a and the rotor. On the other hand, the torque acting between the stator yoke unit b and the rotor is represented as ##EQU2## because the stator yoke unit b deviates in phase by 1/4 pitch relative to the stator yoke unit a. On the other hand, in terms of electrical angle, That is, in contrast with Ta, the sign has only changed. The curve Tb of FIG. 17 of the accompanying drawings represents the above equation with p as p =1.\nHere, the cogging torque To as the stepping motor can be represented by the sum of cogging torques acting on the two stator yoke units a and b and the rotor and therefore is ##EQU3## The curve To of FIG. 16 shows the characteristic of the cogging torque when a.sub.1 &gt;b.sub.1, and the cogging torque To of the stepping motor does not differ from the cogging torque acting between a single stator yoke unit and the rotor, in the number and positions of the points at which the torque is zero, and the peak value thereof only becomes small. If a.sub.1 =b.sub.1, the cogging torque becomes zero. Also, if a.sub.1 and b.sub.1 can be controlled arbitrarily, the magnitude of the cogging torque can be set to any magnitude.\nConsider now a case where deviation exists in the circumferential direction with respect to the rotor shaft 12, that is, a case where the stator yoke unit a and the stator yoke unit b are disposed with phase deviation of 1/4 pitch and the stator yoke unit b deviates from this position by an angle error .alpha.k (electrical angle).\nAssuming that the position of the rotor when the point at which the center of the magnetic pole of the rotor and the center of the magnetic pole of the stator yoke unit coincide with each other is regarded as the reference angle .theta. (mechanical angle), and assuming that the magnitude of the torque produced in the stator yoke unit a is Tk and the magnitude of the torque produced in the stator yoke unit b is Tk+.DELTA.T, the cogging torque T1 acting on the two stator yoke units a and b and the rotor is represented as ##EQU4## and therefore, if it is resolved and put in order, the following is obtained: ##EQU5## But high-order components are not taken into consideration.\nUsually, the cogging torque refers to the maximum value of the above equation and therefore, T1Max=-.sqroot..sub.(2Tksin.alpha.kcos.alpha.k) 2.sub.+(2TKsin.alpha.k+.DELTA.T) 2 is the cogging torque. Accordingly, the cogging torque has its magnitude determined by the angle error .alpha.k and the non-uniformity of the magnitude of the cogging torque between the stator yoke units.\nHere, assuming that .DELTA.T =0 and examining the influence of only the angle error .alpha.k upon the production of the cogging torque, the maximum value of the cogging torque is T1Max=2Tksin.alpha.k and the influence is as shown in the table of FIG. 18 of the accompanying drawings In calculation, one period of the electrical angle has been calculated with respect to a mechanical angle of 30.degree., i.e., a case where the stator has twelve poles, and assuming that the deviation of the mechanical angle is .theta.k, 2k =12.theta.k. From this result, it is seen that if the two sets of stator yoke units are in a right position, no cogging will occur, but is the stator yoke units deviate by 0.5.degree., there will occur cogging corresponding to 20% of the magnitude of the cogging torque produced by one set of stators.\nDescription will now be made of the relation between the stationary position of the rotor and the position in which the cogging torque To is zero.\nGenerally, between the magnetic energy W by the rotor magnet of a stepping motor and the cogging torque, there is the following relation: ##EQU6## Thus, the magnetic energy W can be represented by ##EQU7##\nFIG. 17 shows the relation between the magnetic energy of the above equation and the rotor position. Here, the rotor 13 tries to come to a standstill at a location whereat the magnetic energy is smallest and therefore, in the position of .alpha.=.pi./2, 3.pi./2 . . . . (2.sub.n-1).pi./2, the cogging torque becomes zero and thus, originally, the rotor tries to come to a standstill in this position, but since the magnetic energy W is high, the rotor becomes unstable and does not come to a standstill (the time when the rotor stably comes to a standstill is the time when frictional force is great). After all, it is in the position of .alpha.=0, .pi., 2.pi., . . . , n.pi. that the rotor stably comes to a standstill.\nThat is, the rotor stably comes to a standstill in a position in which the cogging torque curve changes from the negative to the positive, and the rotor does not stably come to a standstill in a position in which the cogging torque changes from the positive to the negative.\nAs is apparent from the foregoing description, the stationary position of the rotor of the stepping motor in a state in which the coil is not excited is .alpha.=0, .pi., 2.pi., . . . , n.pi. (n being an integer).\nThe operation when the rotor rotates by an amount corresponding to one pitch of the stator poles when the coils of the stepping motor shown in FIG. 12 are excited as shown in FIG. 13 will now be described with reference to the developed view of FIG. 14a showing the stator poles and rotor magnetic poles of the stepping motor.\nAs already described, when the coils are not excited, the rotor 13 is in the position of .alpha.=0. When the phase A and phase B are subsequently excited, the magnetic poles 13a come to the position of 1/8 pitch intermediate of the stator poles 1a, 1a, i.e., the position of .alpha.=.pi./4. In the next phase B and phase A, the magnetic poles come to the position of 3/8 pitch intermediate of 1/2 pitch and 1/4 pitch (.alpha.=3.pi./4). In the next phase A and phase B, the magnetic poles come to the position of 5/8 pitch intermediate of 1/2 pitch and 3/4 pitch (.alpha.=5.pi./4). In the next phase B and phase A, the magnetic poles come to the position of 7/8 pitch intermediate of 3/4 pitch and 1 pitch (.alpha.=7.pi./4). Subsequently, in the phase A and phase B, the same thing as the beginning is repeated.\nIn a stepping motor wherein two yoke members each provided with a plurality of pole teeth parallel to a rotor shaft on the inner surface thereof are combined into one pair and which has two pairs of yokes, welding and fixing each pair of yokes when assembling the yokes is described and shown in Japanese Laid-Open Patent Application No. 59-53079.\nThe stationary angle error of the rotor in the excited condition of the coils of the stepping motor will now be described with reference to FIG. 14b. In this figure, the solid line shows the cogging torque To =(a.sub.1 -b.sub.1) sin 2.alpha., and the plus side which is the upper half of the torque of the vertical axis is a torque for returning the rotor with respect to the direction of movement of the rotor, and the minus side which is the lower half is a torque for advancing the rotor.\nIn the non-excited condition, the rotor is stationary in the position of .alpha.=0.\nNow, since the first step of excitation excites the phases A and B, the exciting torque of dotted line 19 acts and tries to advance the rotor 13 to the position of .alpha.=.pi./4 in which the exciting torque is zero. On the other hand, in the position of .alpha.=.pi./4, a cogging torque which is a.sub.1 -b.sub.1 is acting. This cogging torque acts in a direction to return the rotor to the position of .alpha.=0 and therefore, after all, the rotor comes to a standstill in a position wherein the torque which returns the rotor by the cogging torque and the torque which advances the rotor by the torque exciting the phases A and B are balanced with each other. This position is the position of .alpha.=.alpha.1 in FIG. 14b, and is on this side of the position of .alpha.=.pi./4 which is the stationary position of the rotor which should originally be obtained. At the second step, the phases A and B are excited, and the exciting torque of dotted line 20 acts and tries to advance the rotor to the position of .alpha.=3.pi./4 in which the exciting torque is zero.\nOn the other hand, in the position of .alpha.=3.pi./4, a cogging torque which is -(a.sub.1 -b.sub.1) acts. This cogging torque acts in a direction to advance the rotor and therefore, after all, the rotor comes to a standstill in a position wherein the torque which advances the rotor by the cogging torque and the torque which returns the rotor by the torque exciting the phases A and B are balanced with each other. This position is the position of .alpha.=.alpha.2 in FIG. 14b, and is advanced with respect to the position of .alpha.=3.pi./4 which is the stationary position of the rotor which should originally be obtained.\nAt the next third steps, the same thing as the first step happens and the rotor comes to a standstill in a position .alpha.3 which is on this side of the position of .alpha.=5.pi./4 which is the original stationary position of the rotor. Further, at the fourth step, the same thing as the second step happens and the rotor comes to a standstill in a position .alpha.4 advanced with respect to the position of .alpha.=7.pi./4 which is the original stationary position of the rotor. If the angle error of the rotor relative to the position which should originally be obtained is plotted on the vertical axis and the number of steps of the stepping motor is plotted on the horizontal axis, there is obtained the characteristic of FIG. 14c of the accompanying drawings. In the case of the two-phase excitation system, usually the angle error of the rotor has its sign reversed at each step. That is, the amount of angular movement of the rotor changes to small, great and small at each step, and this has led to the important disadvantage that a uniform accurate amount of movement is not obtained.\nAlso, in the rotating condition of the motor, the fluctuation of the cogging torque can be regarded as the fluctuation of an extraneous load applied to the motor, and this has also led to the disadvantage that when an attempt is made to utilize the stepping motor as a multipole brushless DC motor, it is difficult to drive the motor smoothly so that there occurs no fluctuation of the number of rotations, by speed control in a state in which these load fluctuations corresponding to one half of the number of excitation change-overs exist during one full rotation.\nAs regards a motor of small cogging torque, there has been the disadvantage that where the motor is used for positioning or the like, when the driving current of the motor is cut off, the retained torque is small and therefore the set position is liable to deviate due to external factors."}
{"text": "The present invention relates generally to an air pump, and more particularly to a bicycle air pump.\nThe conventional air pump is provided with a large cylinder and a small cylinder, which are shifted by a switching device in the midst of the inflation process. The operation of the switching device is done manually. The manual operation of the switching device causes a great deal of inconvenience to the user of the air pump.\nThe primary objective of the present invention is to provide an air pump which is free from the drawback of the conventional air pump described above.\nThe air pump of the present invention comprises at least two cylinders, a distribution guide tube, and a spring. The distribution guide tube is disposed at the lower end of an outer tube in which the large and the small cylinders are disposed. The distribution guide tube is provided therein with two air chambers, one of which is in communication with the large cylinder and is provided with a pressure leaking slide block in intimate contact with the inner wall of the air chamber. The pressure leaking slide block is provided with a channel which is opened or closed along with the movement of the slide block. The spring is disposed in the air chamber for providing the pressure leaking slide block with the recovery spring force. Another air chamber is in communication with the large and the small cylinders for inflating an inflatable object."}
{"text": "The background description provided herein is for the purpose of generally presenting the context of the disclosure. Work of the presently named inventors, to the extent it is described in this background section, as well as aspects of the description that may not otherwise qualify as prior art at the time of filing, are neither expressly nor impliedly admitted as prior art against the present disclosure. Unless otherwise indicated herein, the approaches described in this section are not prior art to the claims in the present disclosure and are not admitted to be prior art by inclusion in this section.\nDevice-to-device (D2D) communication is a study item for the third generation partnership project (3GPP) long term evolution (LTE) standards. Specifically, inclusion of support for direct communication between two or more devices is being considered for the 3GPP LTE standards. However, in some cases D2D communications and cellular communications, which may involve longer range communication between a user equipment (UE) and eNodeB (eNB) may interfere with one another. For example, if two devices are communicating via D2D communications, then a cellular communication between a UE and an eNB may suffer from interference caused by the D2D communications. Similarly, the communication between an eNB and a UE may cause interference to the D2D communications."}
{"text": "Ion sources used in mass spectrometry produce continuous or quasi-continuous beams of charged particles. Even in the case of pulsed operation of an ion source, accumulation of charged particles during several cycles of operation in a special storage device may be necessary. Therefore, in the case of pulsed operation of mass-analysers, special devices are used to ensure decomposition or breaking-up of a continuous beam of charged particles or the contents of a storage device, into separate portions and transportation thereof to the mass-analyser input. In recent devices used for transportation of charged particles, the tasks of cooling and spatial compression of charged particle packets for the purpose of a reduction of their emittance (the size of a packet of particles in phase-space coordinates) can also be solved efficiently, and additional manipulations can be performed with the charged particles during transportation (for example, fragmentation of charged particles, generation of secondary charged particles, selective extraction of charged particles to be subject to detailed analysis, etc.).\nSeveral types of radio-frequency (RF) devices are used in mass spectrometry for charged particle manipulation. The first group of such devices includes mass analysers (as well as mass separators and mass filters). The purpose of such devices is the selection of those particles featuring particular mass-to-charge ratio, from the totality of charged particles. The main types of RF mass analysers include quadrupole mass filters and ion traps.\nRadio-frequency quadrupole mass filters and ion traps proposed by Paul are known starting from about 1960s. Both types of mass analysers have been proposed in U.S. Pat. No. 2,939,952. Rather recently, linear ion traps were proposed, with radial ejection of charged particles from the trap (U.S. Pat. No. 5,420,425) and ejection of ions from the trap along the axis (U.S. Pat. No. 6,177,680). A detailed description of the principle of operation of said devices can be found, for example, in R. E. March, J. F. J. Todd, Quadrupole Ion Trap Mass Spectrometry, 2nd edition, Wiley-Interscience, 2005; F. J. Major, V. N. Gheorghe, G. Werth, Charged Particle Traps, Springer, 2005; G. Werth, V. N. Gheorghe, F. J. Major, Charged Particle Traps II, Springer, 2009.\nFunctioning of quadrupole mass filters is based on the theory of solution stability of the Mathieu equation (see, for example, N. W. McLachlan, Theory and Application of Mathieu Functions, Claredon Press, Oxford, 1947 (chapter 4) or M. Abramovitz and I. Stegun, Handbook of Mathematical Functions with Formulas, Graphs and Mathematical Tables, 10ed., NBS, 1972 (chapter 20).). In the case of well-selected parameters of the intensity of quadrupole DC electric field, intensity of quadrupole RF field and the frequency of quadrupole RF field, charged particles having a particular mass-to-charge ratio would pass through the RF quadrupole mass filter. The other charged particles would lose the stability of their trajectories, and would be lost outside the boundaries of the channel of the mass filter.\nOperation of mass analysers of the ion trap type is generally based on the theory of the Mathieu equation. In these mass analysers, a quadratic or nearly quadratic electric field is used, obtained through application of ideal hyperbolic electrodes, and the analysers are filled with a light gas at low enough pressure. In such devices, after slowing down the speed of motion of the charged particles due to multiple collisions with the molecules of neutral gas, the particles would then sequentially be extracted from the device by means of swinging/oscillating of the group of charged particles having the required mass-to-charge ratio, with the help of an RF electric field having the required frequency. The picture described above is somewhat approximate, since the practical ion trap mass spectrometry has developed and employed rather sophisticated methods for isolation, fragmentation and selective ejection of charged particles from ion traps by means of the action of specially configured RF fields on the particles.\nAnother important group of RF devices includes RF transporting devices for ion beams. The purpose of such devices is the confining of a beam of charged particles having different masses, within a bounded region inside the device (for example, near the axis of the device), and transfer of charged particles from one point within the space (point of inlet) to another point within the space (point of outlet).\nA wide class of such devices is based on application of a two-dimensional multipole field, or approximate multipole field, extended along the third coordinate. The devices are used, for example, for transfer of ions from gas-filled ion sources operating at rather high gas pressures, into devices for mass-analysis of ions, operating at considerably lower pressure of gas, or in vacuum. Because of the fact that said linear multipole ion traps are not used directly for mass analysis, the requirements towards a strictly quadratic or strictly multipole field would not be vital, and for the purpose of simplification of the production technology while manufacturing such devices, hyperbolic and multipole electrodes, as a rule, would be replaced with cylindrical rods or even more coarsely shaped electrodes.\nWhen charged particles are transferred into a linear multipole trap, collisions of the charged particles with gas molecules reduce their kinetic energy and force the particles to be groped near the axis of the device (U.S. Pat. No. 4,963,736). This ensures such an important function like beam cooling and spatial compressing of the beam of charged particle for the purpose of reduction of the beam emittance (i.e., the volume of an ensemble of charged particles, corresponding to the beam, in phase space). An RF electric field is capable of confining charged particles in a radial direction, at a stage where the reduction of kinetic energy of charged particles has not yet taken place, even in the case of relatively high kinetic energies, and “compresses” the particles towards the axis in the course of the loss of their kinetic energy.\nThe gas-filled linear multipole ion beam transporting devices described above are frequently used simultaneously, as collision cells for fragmentation of charged particles in tandem mass spectrometers (for example, see. U.S. Pat. No. 6,093,929). A DC electric field directed along the axis of the device, the field created by additional electrodes, can be used for forced transfer of charged particles along the channel of transfer (ion transporting device proposed in U.S. Pat. No. 5,847,386, collision cell for fragmentation of ions proposed in U.S. Pat. No. 6,111,250).\nIf the ends of a linear multipole ion transporting device are closed using barriers formed by an electric field, another type of RF device used in mass spectrometry is formed—a linear multipole ion trap, or a storage device for charged particles. Such traps are widely used to accumulate charged particles and pulse transmission of charged particles into an analysing device (U.S. Pat. No. 5,179,278, WO02078046, U.S. Pat. No. 5,763,878, U.S. Pat. No. 6,020,586, U.S. Pat. No. 6,507,019 and GB2388248). Multipole ion traps are also frequently used to initiate task-oriented ion-molecular reactions between charged particles and neutral particles (U.S. Pat. No. 6,140,638 and U.S. Pat. No. 6,011,259), or electrons (patent Nos. GB2372877, GB2403845 and GB2403590), or charged particles with opposite charges (U.S. Pat. No. 6,627,875), to provide additional fragmentation of charged particles due to exposure of the same to an impact, for example, of photons, or other external physical factors.\nThe RF ion trap proposed by Paul, or a linear trap, can also be used for the same purpose as a multipole linear trap, when the total amount of ions is injected at once from the trap into an analysing device due to a pulse of electric voltage, instead of consecutive resonance ejection of the desired groups of ions (patent Nos. WO2006/129068 and US2008/0035841). In a similar way, a multipole linear trap, wherein the injection into the analysing device is made mass-selective, can be used as a rough mass filter, which selects the required groups of charged particles for further detailed analysis (patent No. US2007/0158545).\nThere are devices known to have functions similar to the above-mentioned transporting devices, which include transporting devices and/or storage devices wherein electrodes are used, in the form of an array of plates with apertures, and to which electrodes RF voltages are applied, with phase shift between adjacent plates (U.S. Pat. No. 6,812,453, U.S. Pat. No. 6,894,286 and U.S. Pat. No. 5,818,055), or between the parts forming one plate (patent No. PCT/GB2010/001076). In that case, because of the symmetry of electrodes, the generated RF field near the axis of the device would be practically zero, whereas it would grow abruptly near the boundaries of the transporting channel. Therefore, like in the case of the linear multipole ion transporting devices, the charged particles would be repelled from the electrodes and confined by the RF field within a limited space surrounding the axis of the device, and in the course of reduction of their kinetic energy due to collisions with gas molecules, the charged particles would be grouped near the axis of the device.\nOne can see that in the case of an absence of additional electric fields in the vicinity of the axis of the device, the forces enabling the movement of charged particles along the axis of the transporting device would practically be absent due to symmetry of the electrodes and high frequency of the electric field (U.S. Pat. No. 5,818,055 and U.S. Pat. No. 6,894,286), and the transfer of charged particles along the length of the channel for transportation would not be very efficient. Indeed, the capture of charged particles moving along the axis of the device is not mentioned in U.S. Pat. No. 5,818,055 and U.S. Pat. No. 6,894,286; furthermore, the particles having different masses and different initial conditions (coordinates and velocities) move along the channel of transportation with different effective velocities, and as a result, there would be no separation of the beam of charged particles into individual spatially separated and synchronically transferred packets of charged particles.\nThe superposition of radially non-uniform RF electric field, which enables localisation of charged particles in the vicinity of the axis of the device along the radial direction, and quasi-static progressive wave of electric field along the axis of the device enabling splitting of the beam of charged particles having different masses into spatially separated packets and synchronous transportation of said packets along the axis of the device may be the most successful solution from among the above-mentioned solutions (U.S. Pat. No. 6,812,453 and PCT/GB2010/001076).\nHowever, since the positively charged particles are grouped in the vicinities of minima of the progressive wave of potential of the quasi-static electric field, and negatively charged particles are grouped in the vicinities of maxima of the progressive wave of potential of the quasi-static electric field, it would not be possible to ensure transportation of positively and negatively charged particles in an integrated packet of charged particles using this method.\nThe functioning of the majority of RF mass-spectrometry devices is based on the property of an RF electric field to “eject” the charged particles, regardless of the polarity of their charge, from the area of high amplitude of electric field into the area with lower amplitude of electric field. This property has been the consequence of the inertia of motion of charged particles having non-zero masses, under the influence of a fast oscillating electric field.\nThis phenomena is described quantitatively with the help of the theory of effective potential or pseudopotential, first introduced by P. L. Kapitza (see L. D. Landau, E. M. Lifshitz, Mechanics, Ser. Theoretical Physics, M., Fizmatlit, 2004, p. 124-127; G. M. Zaslavsky and R. Z. Sagdeev, Introduction to nonlinear physics: from pendulum to turbulence and chaos, M., Nauka, 1988, p. 49-51 and p. 52-54; M. I. Yavor, Optics of Charged Particle Analysers, Ser. Advances of Imaging and Electron Physics, Vol. 157, Elsevier, 2009, p. 142-144). That is, suppose the frequency ω of oscillations of electric field {right arrow over (E)}(x,y,z,t), which follows the law {right arrow over (E)}(x,y,z,t)={right arrow over (E)}0(x,y,z)cos(ωt+Φ), is high enough (where {right arrow over (E)}0(x,y,z) is the amplitude of oscillations of electric field in a point within the space (x,y,z), ω—frequency of oscillations, φ—initial phase of oscillations, t—time), and the displacement of charged particle having the mass m and charge q, during one period of oscillations of the electric field is small, then the motion of the charged particle can be represented as an “averaged” or “slow” motion, with an added rapid oscillating motion, featuring, however, small amplitude. In that case, the equation for averaged motion would look like as if the averaged motion takes place within electric field having the potential Ū(x,y,z)=q|{right arrow over (E)}0(x,y,z)|2/(4mω2), where the values q, {right arrow over (E)}0(x,y,z), m and ω characterizing the oscillating electric field and the charged particle, have been defined above. The details and substantiation of the theory can be found in the references cited above.\nDue to the fact that the expression for potential Ū(x,y,z) includes charge q and mass m, the potential Ū(x,y,z) affects equally both positively and negatively charged particles, and the effect is also dependent on the mass of a charged particle. In case of a real electric potential U(x,y,z) positively charged particles would undergo a force directed reversely with respect to the gradient of electrical potential, and negatively charged particles would undergo a force directed along the gradient of electrical potential, whereas such force would not be dependent on the mass of a particle. From the expression for potential Ū(x,y,z) it follows, that a charged particle would be <<pushed out>> from the area where the amplitude of oscillations of the RF field is high, into the area where said amplitude of oscillations of the RF field is lower (that is, from the area where the potential Ū(x,y,z) has a higher value, the particle would move into the area where the potential Ū(x,y,z) has a lower value). The extracting action of the RF electric field is not dependent on the polarity of charged particle, and moves both positive and negative charged particles in the same direction. The extracting action of the RF electric field is weaker with respect to those charged particles having heavier masses, than with respect to lighter charged particles. The extracting action of the RF electric field can be controlled by varying the frequency of oscillations of the electric field.\nThe potential Ū(x,y,z) is called an effective potential, or a pseudopotential, and represents a useful mathematical tool for describing and analysing the averaged motion of a charged particle (though in fact, it does not actually correspond to any physical fields). We shall take for granted, some of its properties. For electric field {right arrow over (E)}(x,y,z,t), which varies with time t under the law of harmonic oscillations {right arrow over (E)}(x,y,z,t)={right arrow over (E)}0(x,y,z)cos(ωt+φ) with a constant amplitude {right arrow over (E)}0(x,y,z) at a point (x,y,z), with a constant frequency ω and with a constant phase shift φ=const, the pseudopotential Ū(x,y,z), which affects a charged particle having the charge q and mass m, is calculated using the above formula Ū(x,y,z)=q|{right arrow over (E)}0(x,y,z)|2/(4mω2). If the phase of the RF field is not constant over the entire space, but varies from point to point in a predetermined manner φ=φ(x,y,z), so that the law of variation of the RF electrical field with time t has a more sophisticated form {right arrow over (E)}(x,y,z,t)={right arrow over (E)}0(x,y,z)·cos(ωt+φ(x,y,z))={right arrow over (E)}c(x,y,z)·cos ωt+{right arrow over (E)}s(x,y,z)·sin ωt, where {right arrow over (E)}c(x,y,z) is the amplitude of harmonic component cos ωt in the point of space (x,y,z), {right arrow over (E)}s(x,y,z) is the amplitude of harmonic component sin ωt in the point of space (x,y,z), and the values {right arrow over (E)}0(x,y,z), ω and φ(x,y,z) were defined earlier, then the pseudopotential Ū(x,y,z) corresponding to the given RF electrical field would be calculated using the formula Ū(x,y,z)=q(|{right arrow over (E)}c|2+|{right arrow over (E)}s|2)/(4mω2), where q is the charge of a particle, and m is its mass. If the RF field under consideration is a time-dependent periodic function, so that the electric filed intensity {right arrow over (E)}(x,y,z,t) in the point of space (x,y,z) at the point of time t can be represented as a Fourier series in the form of {right arrow over (E)}(x,y,z,t)=Σ{right arrow over (E)}c(k)(x,y,z)cos(kωt)+{right arrow over (E)}s(k)(x,y,z)sin (kωt), where {right arrow over (E)}c(k)(x,y,z) is the amplitude of harmonic component cos kωt of electric field in the point of space (x,y,z), {right arrow over (E)}s(k)(x,y,z) is the amplitude of harmonic component sin kωt of electric field in the point of space (x,y,z), k is the number of harmonic component, ω is fundamental frequency of the RF electric field, then the pseudopotential Ū(x,y,z) of such RF electric field would be calculated as a sum of contributions of individual harmonic components, using the formula Ū(x,y,z)=qΣ(|{right arrow over (E)}c(k)(x,y,z)|2+|{right arrow over (E)}s(k)(x,y,z)|2)/(4mω2k2), where q is the charge of a particle, and m is its mass. If in addition to the RF electric field {right arrow over (E)}(x,y,z,t), there is an electrostatic field having potential of U(x,y,z), the electrostatic potential U(x,y,z) and the pseudopotential Ū(x,y,z) would be summed. If there are several different RF electric fields with essentially different frequencies, then individual pseudopotentials would be summed for these electric fields, however, if the difference between the frequencies of these RF fields is insignificant, this rule would not be valid. If, for the purpose of simulation of charged particle kinetic energy reduction as a result of collisions with gas molecules, an effective viscous friction is introduced, having an impact on the charged particle with a force {right arrow over (F)}=−γ({right arrow over (v)}−{right arrow over (v)}0), where {right arrow over (v)}(t) is the velocity of particle at time t, {right arrow over (v)}0(x,y,z) is the velocity of gas molecules in the point (x,y,z), and γ is the viscous friction coefficient, which does not depend on time, coordinates, and electric field, then the result of “slow” motion of charged particle would be as if all the three factors (electrostatic potential, pseudopotential and viscous friction) were affecting the charged particle simultaneously and independently.\nIt should be emphasised that the description of motion of a charged particle, using pseudopotential, only represents a mathematical approximation, obtained under certain assumptions as regards the motion of charged particle, and may not correspond to its actual motion. In this respect, for the purpose of analysis of charged particle motion in the above mentioned radio-frequency quadrupole mass filters and radio-frequency ion traps, it would be necessary to perform a rigorous analysis of motion of a charged particle in the actual electric fields (i.e., Mathieu equation theory), in order to obtain the correct structure of the zones of stability of motion. The approach based on the use of pseudopotential would not give a correct solution, because under the conditions where a charged particle moves near the boundary of the zone of stability, and a resonance takes place between “slow” oscillations of the charged particle and the RF electric field, the displacement of the charged particle during one period of the RF electric field under no conditions could be considered to be small.\nThe present inventors have considered the operation of the device of U.S. Pat. No. 6,812,453 in more detail.\nThe device under consideration contains a system of electrodes representing a series of coaxially positioned plates with apertures arranged to create internal space between the electrodes, the space directed along the longitudinal axis of the device, and intended for transmission of ions within the same. The device also includes a source of power supply, which provides supply voltage to be applied to the electrodes, including alternating high frequency voltage component, the positive and negative phases of which are applied alternately to the electrodes, and quasi-static voltage component, for creation of which, static or quasi-static voltages are applied to the electrodes successively and alternately, in particular, in the form of unipolar or bipolar pulses of a DC voltage.\nThe said device creates an electric field, the intensity of which {right arrow over (E)}(x,y,z,t) is described by the expression {right arrow over (E)}(x,y,z,t)={right arrow over (E)}a(x,y,z,t)+{right arrow over (E)}0(x,y,z)ƒ(t), where {right arrow over (E)}a(x,y,z,t) is a quasi-static electric field varying along the length of the channel for charged particles transportation, depending on the spatial coordinates (x,y,z) and time t, {right arrow over (E)}0(x,y,z) is time-independent and non-uniform, at least in a radial direction, amplitude of the RF electric field, depending on spatial coordinates (x,y,z) and independent on time t, ƒ(t)=cos(ωt+φ) is the rapidly oscillating function of time t, which in this particular case describes strictly harmonic oscillations with the frequency ω and initial phase φ. Quasi-static behaviour of the function {right arrow over (E)}a(x,y,z,t) and the rapidness of oscillations of the function ƒ(t) are understood in the sense that during a period where the function ƒ(t) has time to perform several oscillations, the function {right arrow over (E)}a(x,y,z,t) remains practically unchanged. Mathematical notation of this condition is written in the form of inequality |∂{right arrow over (E)}a/∂t|2/|{right arrow over (E)}0|2<<|df/dt|2, which should be satisfied, in order that the device would function properly. Thereby variation of the electric field {right arrow over (E)}(x,y,z,t) with time would have two time scales: a “fast time”, during which the value of the function {right arrow over (E)}0(x,y,z) ƒ(t) would be noticeably changed, and a “slow time”, during which the value of the function {right arrow over (E)}a(x,y,z,t) would be noticeably changed.\nFIGS. 1 to 9 assist with understanding the operation of the device of U.S. Pat. No. 6,812,453. FIG. 1 demonstrates a round diaphragm used as a single electrode for the device according to U.S. Pat. No. 6,812,453. FIG. 2 shows the arrangement of the aggregate of round diaphragms with respect to the channel for charged particles transfer, according to U.S. Pat. No. 6,812,453. FIG. 3 shows the distribution of axial component of the intensity of electric field according to U.S. Pat. No. 6,812,453 along the length of the channel for charged particle transportation, for a series of close points in time t, t+δt, t+2δt, t+3δt, . . . (that is, in a “fast” time scale). FIG. 4 shows variation of the envelope of axial component of the electric field of U.S. Pat. No. 6,812,453 along the length of channel, for a number of points in time t and t+Δt, located sufficiently far from each other (that is, in a “slow” time scale). The radial component of the electric field equals zero at the axis of the device of U.S. Pat. No. 6,812,453 due to the symmetrical configuration of the electrodes. FIG. 5 shows a two-dimensional distribution of pseudopotential Ū0(x,y,z) along the length of the channel for charged particle transportation, and in a radial direction of the channel for transportation, which corresponds to the RF electric field according to U.S. Pat. No. 6,812,453. FIG. 6 shows possible two-dimensional distribution (at some point in time) of the potential Ua(x,y,z,t) of the quasi-static electric field {right arrow over (E)}a(x,y,z,t) of U.S. Pat. No. 6,812,453. FIG. 7 shows possible distribution of the potential Ua(x,y,z,t) of quasi-static electric field {right arrow over (E)}a(x,y,z,t) of U.S. Pat. No. 6,812,453, along the length of the channel for charged particle transportation. FIG. 8 shows possible summary electric voltages, which can be applied to the first, second, third, fourth electrode, respectively, in each group of four repetitive electrodes, according to U.S. Pat. No. 6,812,453. (In these examples, the simplest possible case is considered, of the progressive wave of quasi-static potential Ua(x,y,z,t), formed along the channel intended for the motion of charged particles, according to U.S. Pat. No. 6,812,453, viz., the case of a wave having purely sinusoidal waveform.)\nAccording to U.S. Pat. No. 6,812,453 the charged particles are “forced” towards the axis of the device as a result of the action of the RF field and formation of the pseudopotential Ū0(x,y,z) over the radius thereby forming a barrier farther from the axis of the device, and after damping of kinetic energy to equilibrium value, appear to be collected in the neighbourhood of the axis of the device. Due to the presence of the distribution of the quasi-static electric potential with alternating local minima and maxima along the axis of the device, positively charged particles are not just concentrated around the axis of the device, but are collected in local minima of the quasi-static electric potential, as soon as their kinetic energy proves to be lower than the local maxima of the quasi-static electric potential. Respectively, the negatively charged particles, after cooling as a result of collisions with gas molecules, are collected in local maxima of the quasi-static electric potential (the positively charged particles are affected by the force directed against the gradient of the electric potential, while negatively charged particles are affected by the force directed along the gradient of the electric potential).\nThe fact that at some interval along the length of the axis (in particular, in the neighbourhood of the minima of electric potential for positively charged particles and in the neighbourhood of the maxima of electric potential for negatively charged particles), while moving away from the axis, the radial electric field of quasi-static potential repels the charged particles from the axis of the device, is of no importance, since the repelling action of the RF field, returning the charged particles back to the axis of the device is overbalancing i.e. dominant. When the wave of the quasi-static potential Ua(x,y,z,t) travels slowly along the axis of the device, it captures the charged particles, located near the axis of the device in the neighbourhood of local maxima and minima of the quasi-static potential, while forcing the particles having different masses and different kinetic energies to move synchronously. The process is shown schematically in FIG. 9. Note that this results in alternating groups of positively and negatively charged particles.\nNumerical simulation by the present inventors of the actual motion of charged particles in the described electric fields confirms this qualitative picture of motion. For output devices operating in pulsed mode, this method of separation of a continuous flow of charged particles into discrete portions seems to be the most successful. With a correct setting of time intervals between arrivals of individual discrete portions of charged particles from the output of the transporting device and correspondingly, to the input of the next device (which, as a rule, represents a mass analyser operating in pulsed mode), and the time of the next analysis of arrived portion of charged particles, this method allows analysis of all the charged particles from the continuous beam into the analyser, practically without losses.\nHowever, the device of U.S. Pat. No. 6,812,453 does not provide a capability of combining positively and negatively charged particles in a single transported packet."}
{"text": "In the machining industry, when a piece of stock or a workpiece is being machined, it is desirable to securely hold the workpiece so as to prevent unwanted movement thereof during the machining operation. One approach commonly taken is to utilize any one of a variety of vise type devices for holding workpieces. Such vise devices typically include one or more jaws which can be adjusted to engage the outer surfaces of the workpiece. Often in the machining environment, it is necessary to utilize the same machine tool for machining workpieces of various configurations, that is, workpieces having different cross-sectional shapes. In the past, some vise type devices utilized for holding workpieces have allowed for the jaw or jaws to be removed in order to accommodate holding of workpieces having various configurations. In particular, it is not uncommon to use multiple sets of jaws, each set being machined to specifically hold a workpiece of a certain configuration. One set of jaws is placed on the vise type device at a time in order to hold a workpiece having a configuration which corresponds thereto. If the same vice needs to be used for holding a workpiece of a different configuration the solution is to remove the set of vise jaws currently being used and to replace the removed set with a different set of vise jaws. This type of change out operation is time consuming and therefore costly in that it reduces productivity.\nAn alternative to changing out jaws is to use a jaw which does not necessarily match any one of the workpiece configurations. However, this can sometimes result in inadequate holding forces, and thus workpiece movement during the machining operation, resulting in defective end products. Even where specially machined jaws are used, it is possible that some workpieces may not exactly fit the jaw. For example, a jaw might be machined having planar surfaces intersecting at an angle of 90 degrees in order to hold a workpiece of rectangular cross-section. It is possible that some workpieces may not be exactly rectangular, that is, the corners thereof may not form precise right angles. In such cases the machined jaws are unable to compensate for such variations and the effective holding power of the jaws is adversely affected with respect to such imperfect workpieces.\nFurther, many of the aforementioned problems arise in workpiece holding applications in which a device other than a vice is used. It would therefore be desirable to provide a workpiece holding assembly which can be adapted for inclusion in a large number of types of holding devices, and which can be scaled for use in holding both small and large workpieces.\nAccordingly, the present invention is directed to overcoming one or more of the problems as set forth above."}
{"text": "1. Field of the Invention\nThe present invention relates to a distillation apparatus that carries out a distillation operation widely applied to many industrial processes, and more particularly to a heat integrated distillation apparatus.\n2. Description of the Related Art\nDistillation separation is a unit operation widely applied to industrial processes in general, but consumes a large amount of energy. In the industrial field, therefore, studies have been conducted on an energy saving distillation systems. Such studies have brought about development of a heat integrated distillation column (hereinafter, HIDiC) as a distillation apparatus that save much energy.\nAs shown in FIG. 1, a basic system of the HIDiC has a structure where a rectifying section (high-pressure unit) and a stripping section (low-pressure unit) are provided such that they are separate from each other. Operation pressure of the rectifying section is set higher than that of the stripping section so that the operation temperature of the rectifying section can be higher than that of the stripping section. This enables a reduction in the amount of heat that is supplied to a reboiler because heat transfer occurs from the rectifying section to the stripping section when there is a heat-exchange surface therebetween. Heat of the rectifying section moves to the stripping section, and hence the amount of heat that is supplied at a reboiler can be reduced. As a result, high energy saving distillation apparatus can be achieved.\nIn order to put the concept of HIDiC to practical use, a number of distillation apparatuses having double-pipe structures, that is, double-pipe structures constituted of inner pipes forming rectifying sections and outer pipes forming stripping sections (refer to JP2004-16928A) have been proposed. These configurations are described as being capable of reducing the amounts of heat that are supplied to the reboilers and the amounts of heat that are removed at the condensers, since heat transfer occurs from the rectifying sections (inner pipes) to the stripping sections (outer pipes).\nHowever, the heat integrated distillation apparatus having the rectifying section and the stripping section formed into the double-pipe structures as discussed in Patent Literature 1 had the following problems 1) to 6).\n1) The product cannot be obtained with side-cut stream. The side-cutting means that a product is withdrawn as an intermediate distillate product, during a distillation process until an end distillate is acquired from top of column.\nIn the distillation apparatus described in JP2004-16928A, the tube units of the double-pipe structures are arranged to come into contact with each other. Moreover, the outer pipes and the inner pipes are equipped with the structured packing. As a result, no pipe arrangement can be formed to withdraw any intermediate distillate product from the inner pipe of each tube unit. Consequently, the structure disables side-cutting.\n2) The feed stage where feed stream is provided cannot be optimized. This is because in the rectifying section and the stripping section formed into the double-pipe structures, packing heights thereof are equal, disabling free setting of the number of stages of the rectifying section and the stripping section.\n3) The feed stage cannot be changed so as to meet the feed stream composition. This is because of the structure where free setting of the feeding stage position is disabled as described in 2).\n4) Multi-feed stream (reception of a plurality of feed streams) cannot be dealt with. This is because of the structure where no feed stream can be supplied in the midway of the double-pipes as described in 1).\n5) Maintenance of the apparatus is difficult. The tube units that use the structured packing are densely arranged to be adjacent to each other as described in 1). This disables complete access to the desired tube unit, and maintenance thereof cannot be carried out.\n6) The heat exchanged rate between the rectifying section and the stripping section that uses double-pipes and in which there is no a degree of freedom in design for designing the heat transfer area, depends only on the temperature profile of the distillation column. Hence, in apparatus design, a degree of freedom in design of heat exchanged rate is small.\nQ, the heat exchanged rate between the rectifying section and the stripping section, is represented by Q=U×A×ΔT, where U is an overall heat-transfer coefficient, A is a heat transfer area, and ΔT is a temperature difference between the rectifying section and the stripping section. In the HIDiC using the double-pipe structure, an inner pipe wall surface becomes a heat transfer area. This heat transfer area has a fixed value determined by a structure of the double-pipes. The overall heat-transfer coefficient also has a fixed value determined by the heat transfer structure and fluid physical properties involved in heat exchange. Thus, as can be understood from the heat exchanged rate formula, a heat exchanged rate on design specification can be changed based only on the temperature difference between the rectifying section and the stripping section, which is changed by the operating pressure of the rectifying section and the stripping section.\nAs the heat integrated distillation apparatus that can solve the problem as described above, the present applicant has proposed the apparatus of JP4803470B.\nFIG. 2 shows a first example of the distillation apparatus disclosed in JP4803470B. The distillation apparatus includes rectifying column 1, stripping column 2 located higher than rectifying column 1, first pipe 23 for communicating column top 2c of the stripping column with column bottom 1a of the rectifying column, and compressor 4 configured to compress vapor from column top 2c of the stripping column to feed the compressed vapor to column bottom 1a of the rectifying column. The distillation apparatus further includes tube-bundle-type heat exchanger 8 located at a predetermined stage of rectifying column 1, liquid withdrawal unit 2d located at a predetermined stage of stripping column 2 and configured to withdraw a part of liquid from the predetermined stage to the outside of the column, second pipe 24 for introducing the liquid from liquid withdrawal unit 2d to heat exchanger 8, and third pipe 25 for introducing fluids introduced through second pipe 24 to heat exchanger 8 and then discharged out of heat exchanger 8 to a stage directly below liquid withdrawal unit 2d. \nWith the above described configuration in FIG. 2, the fluids flow from stripping column 2 to heat exchanger 8 of rectifying column 1 through second pipe 24. Heat is removed from the vapor of rectifying column 1 in heat exchanger 8. Then, the heat can be transferred from rectifying column 1 to stripping column 2 through third pipe 25. The fluids flow from stripping column 2 to rectifying column 1 by gravity. The fluids in heat exchanger 8 are accordingly pushed to flow from rectifying column 1 to stripping column 2. In other words, this heat integrated distillation apparatus employs a thermo-siphon system, and hence no pressure-feeding means such as a pump is necessary for supplying the liquid from rectifying column 1 to stripping column 2 located above in a vertical direction.\nFIG. 3 shows a second example of the distillation apparatus disclosed in JP4803470B. The distillation apparatus includes rectifying column 1, stripping column 2 located higher than rectifying column 1, first pipe 23 for connecting column top 2c of the stripping column with column bottom 1a of the rectifying column, and compressor 4 that compresses vapor from column top 2c of the stripping column to feed the compressed vapor to column bottom 1a of the rectifying column. The distillation apparatus further includes liquid sump unit 2e located at a predetermined stage of stripping column 2 and configured to hold liquid that has flowed downward, heat exchanger 8 located in liquid sump unit 2e, partition plate 16 that is set in a predetermined position of rectifying column 1 and configured to apart upper stages and lower stages completely, second pipe 29 for introducing vapor below partition plate 16 to heat exchanger 8, and third pipe 30 for introducing fluids introduced through second pipe 29 to heat exchanger 8 and then discharged out of heat exchanger 8 to an upper side of partition plate 16.\nWith the above described configuration in FIG. 3, the vapor is withdrawn from rectifying column 1 through second pipe 29. The vapor is introduced to heat exchanger 8 in stripping column 2. Then, heat can be transferred from rectifying column 1 to stripping column 2. High-pressure vapor in rectifying column 1 ascends through second pipe 29 to heat exchanger 8 in stripping column 2. A fluid partially or totally condensed from the vapor in heat exchanger 8 is accordingly pushed out from stripping column 2 to third pipe 30 outside the column. Thus, this configuration also necessitates no pressure-feeding means such as a pump in supplying liquid from stripping column 2 to rectifying column 1 located at a lower side in a vertical direction.\nThe apparatus configurations of FIGS. 2 and 3 described above are each capable of reducing the amount of heat that is removed at condenser 7 which is attached to a column top of rectifying column 1, and reducing the amount of heat of reboiler 3 attached to a column bottom of stripping column 2 more, as compared with an ordinary distillation apparatus which has a column in which an upper side is a rectifying section and a lower side is a stripping section with a feed location as a boundary thereof, and which is not of a heat integrated type. As a result, it is possible to provide an energy-efficient distillation apparatus.\nRectifying column 1 and stripping column 2 can be configured by using trayed sections or packed bed sections similar to those of a general distillation apparatus. Hence, the apparatus can deal with side cutting or multi-feed stream without the need for any improvement, and it is possible to easily perform maintenance of the apparatus. For the same reason, the number of stages of the rectifying column or the stripping column can be freely set, and a feed stage can be optimized.\nA heat transfer area can be freely set, and hence the heat exchanged rate can be determined without any dependence on the temperature difference between the columns.\nAs described above, according to the apparatus example described in JP4803470B (FIGS. 2 and 3), energy efficiency is high, side-cutting and setting of a feed stage position can be easily dealt with, and maintenance of the apparatus is easy. Further, the apparatus of the present invention has a structure in which a degree of freedom in design is high, and hence can be easily accepted by the user side.\nIncidentally, concerning the distillation apparatuses shown in FIGS. 2 and 3, the present inventors aim at further enhancement in energy efficiency, and the respective distillation apparatus examples still have the following room to be improved.\nIn other words, in the distillation apparatus shown in FIG. 2, the following method is adopted. Part or all of a liquid in an arbitrary stage of stripping column 2 is removed through pipe 24 outside the column, and is supplied to tube-bundle-type heat exchanger 8 located at an arbitrary stage of rectifying column 1, where heat exchange is performed. Thereafter, a part or all of the amount of the liquid which is removed from stripping column 2 is vaporized by vapor in rectifying column 1 at a higher temperature, and returns to directly below the above described liquid removal position of stripping column 2 via pipe 25 outside the column by a thermo-siphon effect, without energy given from outside by a pump or the like. Such circulation of the fluids is performed.\nIn such a method, a liquid head is needed at the supply side of tube-bundle-type heat exchanger 8 (pipe 24 outside the column) in order to perform circulation of the fluids by the thermo-siphon effect. In other words, as the portions extending in the vertical direction, of pipes 24 and 25 become long correspondingly to the distance (height) between liquid withdrawal position X from stripping column 2 and heat exchanger installation position Y of rectifying column 1, pressure loss through pipe 25 increases. Hence, in order to circulate the fluids by surpassing this, the liquid head based on the inlet position of heat exchanger 8 (end portion of pipe 24 connected with heat exchanger 8) also becomes large. In the tube of heat exchanger 8, however, the pressure becomes high and the boiling point increases due to the increase in the liquid head. Therefore, the temperature difference between the inside of the tube and the outside (shell) of the tube in heat exchanger 8 becomes small correspondingly to the increase of the boiling point. In order to compensate this, a necessity arises to increase the pressure of rectifying column 1, that is, to increase the temperature in rectifying column 1 by increasing the compression ratio of compressor 4. Thus, there is a room to be improved from the viewpoint of energy saving performance.\nIn other words, in the distillation apparatus shown in FIG. 3, the following method is adopted. Partition plate 16 that completely partitions the inside of the column to an upper side and a lower side is installed in an arbitrary stage of rectifying column 1, all of vapor ascending from below partition plate 16 is withdrawn from the column through pipe 29, and is supplied to tube-bundle-type heat exchanger 8 installed at an arbitrary stage of stripping column 2, where heat exchange is performed. Thereafter, a fluid partially or totally condensed in heat exchanger 8 flows through pipe 30 outside the column to the upper side of partition plate 16 in rectifying column 1 by gravity, and the condensed liquid flows through another pipe 31 to be movable to below partition plate 16. Such circulation of the fluids is performed.\nSuch a method intend to withdraw all of the vapor in rectifying column 1 to the outside of the column, and hence adopts a complicated structure in which partition plate 16 is installed in rectifying column 1, and the condensed liquid fed onto partition plate 16 from stripping column 2 is further transferred to a lower side space of partition plate 16 through pipe 31 and control valve 17 outside the column. Thus, there is a room to be improved from the viewpoint of the structure and manufacturing cost.\nFurther, drive force for the fluids passing through the tube of heat exchanger 8 is obtained by giving pressure loss at the upper and lower sides of partition plate 16, and hence pressure of column bottom 1a needs to be made larger than pressure of column top 1c of rectifying section 1 correspondingly to the pressure loss at the upper and lower sides of partition plate 16. Thus, there arises a need for setting pressure to be higher at an outlet side of compressor 4 (namely, increase a compression ratio) correspondingly to increase in the pressure at column bottom 1a side. Therefore, there is also a room to be improved from the viewpoint of energy saving performance.\nIn order to further improve both of the apparatus configuration examples of FIGS. 2 and 3 as above, the present inventors pay attention to the circumstances as follows.\nWhen the stages for performing side heat exchange between the stripping section and the rectifying section are optimally selected, in order to put the concept of HIDiC to practical use, the lowermost stage of the rectifying section is not used for heat exchange with the stripping section in some cases. FIG. 4 shows a conceptual configuration of HIDiC in this case, and FIGS. 5 and 6 show examples of the mode of carrying out HIDiC in this case. In particular, FIG. 5 shows application of the heat transfer system of the apparatus of FIG. 2 to the conceptual configuration of FIG. 4, and FIG. 6 shows application of the heat transfer system of the apparatus of FIG. 3 to the conceptual configuration of FIG. 4.\nAs is understandable with reference to the configurations of FIGS. 4 to 6, a single stage or a plurality of stages in rectifying section lower part 1d is not or are not involved in heat exchange at all in many cases. This not only applies to the tray, but also to the packed bed layer. In the configurations as above, rectifying section lower part 1d which has the highest temperature is not effectively used for heat exchange.\nIn the configurations in which rectifying section lower part 1d is not used for heat exchange with the stripping section as in FIGS. 4 to 6, there is no need to connect the outlet pipe of compressor 4 to column bottom 1a (the lowermost portion of the rectifying section) of the rectifying column. Rather, connecting outlet pipe 4a of compressor 4 to the heat exchanging section (heat exchanging section directly above rectifying section lower part 1d) which is located at the lowest position of the rectifying section does not waste heat. However, if outlet pipe 4a of compressor 4 is directly connected to the heat exchanging section located at the lowermost position of the rectifying section without being directly connected to the lowermost portion of the rectifying section, there will not be any gas in lower part 1d, which is located below this, of the rectifying section, and therefore, distillation operation which has liquid-vapor equilibrium as the principle of separation is not established. Therefore, the apparatus which uses this method of connection cannot be realized.\nThus, the present inventors have decided to move and dispose region 2g (the region shown by the dotted line in FIG. 4), that corresponds to lower part 1d of the rectifying section, to a location that is above the upper part of the stripping column (in other words, above feed stage 2f), as shown in FIG. 4, for example. The disposition like this does not change the flow itself of the fluids at all, and if the stage of lower part 1d of the rectifying section is moved and disposed above feed stage 2f, region 2g, that corresponds to the lower part of the rectifying section, can be operated under the pressure of stripping column 2 which is lower than the pressure of rectifying column 1. As a result, the relative volatility in region 2g becomes large, and the energy (heat amount) itself which is originally necessary for a separating process can be reduced.\nFurther, in the case of the configuration in which outlet pipe 4a of compressor 4 is connected to the position directly above lower part 1d of the rectifying section, as shown by the dotted lines in FIGS. 4 to 6, the vapor from outlet pipe 4a of compressor 4 is supplied into rectifying column 1, and from rectifying column 1, heat is transferred to heat exchanging section 2h in the lower part of stripping column 2, and returns to the position of rectifying column 1 to which outlet pipe 4a is connected, again. The present inventors have considered that, if that is the case, adoption of such an apparatus configuration is more preferable, that outlet pipe 4a of compressor 4 be directly connected to heat exchanging section 2h at the low part of stripping column 2 without being connected all the way to rectifying column 1, and that the fluids which are subjected to heat exchange at heat exchanging section 2h be introduced into rectifying column 1.\nIn the case of the apparatus configuration based on the above concept, the apparatus configuration of FIG. 5 becomes an apparatus configuration of FIG. 7, and the apparatus configuration of FIG. 6 becomes an apparatus configuration of FIG. 8, though the details thereof will be described later. According to the apparatus configurations of FIGS. 7 and 8, among a plurality of heat transfer systems which perform heat exchange between several stages of respective low-pressure column and a high-pressure column, in one heat transfer system provided between the lower part of a low-pressure column and the lower part of a high-pressure column, high-pressure vapor from compressor 4 is directly fed to heat exchanger 8 at the lower part of the low-pressure column, and the high-pressure vapor which gives heat to the lower part of the low-pressure column through heat exchanger 8 is introduced into rectifying column 1, as shown in FIGS. 7 and 8. In this heat transfer system, both the apparatus configuration of FIG. 2 and the apparatus configuration of FIG. 3 are improved. The reason for this is as follows.\nThat is to say, in the apparatus examples of FIGS. 2 and 5, circulation of the fluids by the thermo-siphon effect is used in order to perform heat exchange between the rectifying section and the stripping section, and in order to cause the fluid to circulate, a liquid head is necessary at the supply side (pipe 24 outside the column) of tube-bundle-type heat exchanger 8. However, the apparatus configuration of FIG. 7 does not need a liquid head in some heat transfer systems, and therefore, improvement in energy saving performance is expected.\nMeanwhile, in the apparatus example of FIG. 3, the high-pressure vapor which is supplied to rectifying column 1 moves into stripping column 2 through pipe 29 outside the column, and the fluids that partially or totally condense in heat exchanger 8 of liquid sump unit 2e in stripping column 2 return to rectifying column 1 through another pipe 30. For this purpose, it is necessary to partition the inside of rectifying column 1 completely with partition plate 16, connect pipe 29 to a lower space of partition plate 16, and connect pipe 30 to an upper space of partition plate 16 to set the pressure of the lower space of partition plate 16 to be higher. But, the apparatus configuration of FIG. 8 does not require a pressure difference in some heat transfer systems, and therefore, improvement in energy saving performance is expected."}
{"text": "Traditional camera technology has for many years relied upon the provision of an optical processing system which relies on a negative of an image which is projected onto a photosensitive film which is subsequently chemically processed so as to “fix” the film and to allow for positive prints to be produced which reproduce the original image. Such an image processing technology, although it has become a standard, can be unduly complex, as expensive and difficult technologies are involved in full color processing of images. Recently, digital cameras have become available. These cameras normally rely upon the utilization of a charged coupled device (CCD) to sense a particular image. The camera normally includes storage media for the storage of the sensed scenes in addition to a connector for the transfer of images to a computer device for subsequent manipulation and printing out.\nSuch devices are generally inconvenient in that all images must be stored by the camera and printed out at some later stage. Hence, the camera must have sufficient storage capabilities for the storing of multiple images and, additionally, the user of the camera must have access to a subsequent computer system for the downloading of the images and printing out by a computer printer or the like.\nFurther, digital camera devices have only limited on board processing capabilities which can only perform limited manipulation of sensed image. The main function of the on board processing capability is to store the sensed image. As it may be desirable to carry out extensive modification of an image, the capabilities of such digital camera devices are considered inadequate."}
{"text": "Microelectronic circuits contain various components that produce electromagnetic radiation capable of influencing other circuit components. Switching noise, voltage droop, parasitic inductance, and similar issues can negatively affect circuit performance if left unaddressed. Accordingly, decoupling (or “bypass”) capacitors have been used in microelectronic circuits to decouple one part of an electrical circuit from another so that the negative effect of circuit noise is lessened.\nFor example, a ceramic capacitor has pairs of electrodes separated by dielectric material. The capacitance C of the capacitor is given by the formula C=(kAn)/d, where k is the dielectric constant of the dielectric material, A is the active area of the electrodes, n is the number of layers (with at least one of the layers consisting of an electrode pair plus the separating dielectric), and d is the thickness of each layer. As may quickly be seen from the formula, one way to increase capacitance is to decrease the dielectric layer thickness. The dielectric layer thickness can be reduced to about 1 micron in current-generation multiple-layer ceramic capacitors (MLCCs).\nFor simplicity and clarity of illustration, the drawing figures illustrate the general manner of construction, and descriptions and details of well-known features and techniques may be omitted to avoid unnecessarily obscuring the discussion of the described embodiments of the invention. Additionally, elements in the drawing figures are not necessarily drawn to scale. For example, the dimensions of some of the elements in the figures may be exaggerated relative to other elements to help improve understanding of embodiments of the present invention. The same reference numerals in different figures denote the same elements.\nThe terms “first,” “second,” “third,” “fourth,” and the like in the description and in the claims, if any, are used for distinguishing between similar elements and not necessarily for describing a particular sequential or chronological order. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the invention described herein are, for example, capable of operation in sequences other than those illustrated or otherwise described herein. Similarly, if a method is described herein as comprising a series of steps, the order of such steps as presented herein is not necessarily the only order in which such steps may be performed, and certain of the stated steps may possibly be omitted and/or certain other steps not described herein may possibly be added to the method. Furthermore, the terms “comprise,” “include,” “have,” and any variations thereof, are intended to cover a non-exclusive inclusion, such that a process, method, article, or apparatus that comprises a list of elements is not necessarily limited to those elements, but may include other elements not expressly listed or inherent to such process, method, article, or apparatus.\nThe terms “left,” “right,” “front,” “back,” “top,” “bottom,” “over,” “under,” and the like in the description and in the claims, if any, are used for descriptive purposes and not necessarily for describing permanent relative positions. It is to be understood that the terms so used are interchangeable under appropriate circumstances such that the embodiments of the invention described herein are, for example, capable of operation in other orientations than those illustrated or otherwise described herein. The term “coupled,” as used herein, is defined as directly or indirectly connected in an electrical or non-electrical manner. Objects described herein as being “adjacent to” each other may be in physical contact with each other, in close proximity to each other, or in the same general region or area as each other, as appropriate for the context in which the phrase is used. Occurrences of the phrase “in one embodiment” herein do not necessarily all refer to the same embodiment."}
{"text": "Ever-stricter requirements for the protection of vehicle passengers and other road users place conflicting demands on the developers of motor vehicle bodies. On the one hand, a high dimensional stability of the body is required for the passenger protection, and on the other hand, the body is to also be flexible in order to minimize the risk of injury to road users outside of the vehicle in the event of an accident. In addition, these demands are to be fulfilled with a weight of the body, which is as light as possible, so as to be able to also attain a minimization of the fuel consumption at the same time.\nSheet metal workpieces with locally variable thickness, so-called tailored blanks, which are obtained, e.g. by butt welding sheet metal of different thicknesses or by welding or adhering reinforcements to a substrate sheet metal, is a known technology, which is used in an attempt to meet these demands. A weight minimization is possible here in that the wall thickness of the tailored blanks is adapted exactly to the load distribution in the later finished workpiece and in that highly loaded areas receive a high wall thickness and lightly loaded areas receive a small wall thickness.\nOne disadvantage of this technology is that it requires the production and assembling of various individual parts. In particular, the welding together of the individual parts takes considerable time, because the spots to be welded by a welding tool must be reached and treated one after the other, and because the effort associated with the production and assembling of the individual parts becomes larger, the more exactly the shape and/or number thereof is adapted to the load distribution in the finished workpiece."}
{"text": "1. Technical Field\nThe present invention relates to a semiconductor device and method for manufacturing the semiconductor device.\n2. Related Art\nA typical conventional process for forming a generally-used base layer in a hetero bipolar transistor is illustrated in FIGS. 11A to 11C and FIG. 12. The following description will be made in reference to selectively growing crystal of SiGe for forming the base layer. A silicon oxide film 14, a p-type polycrystalline silicon film 16 and a silicon nitride film are formed on a substrate 10, which is provided with an n-type collector layer 12, and a concave portion is provided in the silicon nitride film and the p-type polycrystalline silicon film 16. Further, the side wall of the concave portion is covered with a silicon nitride film to form a silicon nitride film 18. Then, a portion of the silicon oxide film 14 exposed over the bottom of the concave portion is etched, and further, a side etching in the transverse direction is conducted to provide an opening in the silicon oxide film 14 (FIG. 11A).\nNext, a contamination in the surface is removed within hydrogen atmosphere at higher temperature, and then the base layer is formed in the bottom of the opening of the concave portion via a chemical vapor deposition (CVD) process. Available chemical compounds for a source gas for such CVD process includes hydrogen (H2), dichlorosilane (SiH2Cl2), germane (GeH4), diborane (B2H6), methylsilane (SiH3CH3) and the like. As shown in FIG. 11B, undoped SiGe is grown on the n-type collector layer 12 via an epitaxial process to form an SiGe monocrystalline layer 20. In such case, SiGe is grown on the exposed portion of the p-type polycrystalline silicon film 16 formed in the upper portion of the opening and exposed in the lower surface of a visor section of the p-type polycrystalline silicon film 16 including silicon nitride film 18 to form the SiGe polycrystalline layer 30.\nFurther, a p-type SiGe monocrystalline layer 22 doped with p-type impurity is formed, and SiGe is further deposited on the SiGe polycrystalline layer 30 to form a p-type SiGe polycrystalline layer 32 (FIG. 11C).\nSubsequently, undoped SiGe monocrystalline layer 24 having an altered germanium ratio of SiGe and an undoped silicon layer 26 are respectively grown on the SiGe monocrystalline layer 22 to form a monocrystalline growth layer. In such case, an undoped SiGe polycrystalline layer 33 and an undoped silicon polycrystalline layer 34 are also deposited on the p-type SiGe polycrystalline layer 32 to form a polycrystalline layer, so that the polycrystalline layer is in con-tact with the monocrystalline layer and the polycrystalline layer also provides electrical coupling between the p-type polycrystalline silicon film 16 serving as the base drawing electrode and the base layer (FIG. 12).\nTypical bipolar transistors having a base layer formed in the interior of an opening provided in the insulating layer are described in, for example, Japanese Patent Application Publication No. 2004-193454 and Japanese Patent Application Publication No. 2003-188177. Other related art documents includes Japanese Patent Application Publication No. 2003-45885 and Japanese Patent Application Publication No. 2005-57171.\nNevertheless, as shown in FIG. 12, a polycrystalline layer downwardly growing from the p-type polycrystalline silicon film 16 extends toward the transverse direction after being contacted with the monocrystalline layer, eventually being in direct contact with the portion of the silicon nitride film 18 exposed in the lower surface of the visor section. Larger dimensional area of such contact causes a larger stress generated due to a difference in thermal expansion coefficient between the silicon nitride film and the polycrystalline layer, leading to a creation of recombination energy level in vicinity of the emitter-base junction, so that a stable current gain is difficult to be obtained."}
{"text": "The present invention relates generally to tamper-evident and so-called tamper-resistant closures and, more particularly, to tamper-evident closures for containers whose contents are maintained at an overpressure with respect to atmospheric pressure.\nRecent unfortunate instances of tampering with the contents of containers intended for human consumption have resulted in the design of a variety of container closures provided with tamper-evident features. Such closures attempt to provide readily visible evidence that the closure seal has at some time been disturbed.\nThe present invention is specifically directed to providing tamper-evident closure apparatus for containers whose contents are maintained at an overpressure with respect to atmospheric pressure. An example to which the present invention is especially suited is carbonated beverage containers.\nTamper-evident packaging presently used in conjunction with closures for sealing internally pressurized containers generally fall into one of several types. Shrink seals comprise bands or wrappers, some with distinctive designs, which are shrunk by heat or drying to seal the union of the closure and container. The seal must be cut or torn to open the container. Another type of tamper-evident closure is the breakable closure. The container is sealed by a plastic or metal closure that either breaks away completely when removed from the container or leaves part of the closure attached to the container. The closure must be broken to open the container and remove the product. A metal closure is disclosed in U.S. Pat. No. 3,303,955 issued Feb. 14, 1967 to Osborne, et al. which is provided with a locking band connected to the closure by frangible bridges which fracture to leave the band on the container upon initial removal of the closure from the container to thereby provide visible evidence of initial closure removal. Another metal closure is disclosed in U.S. Pat. No. 3,601,273 issued Aug. 24, 1971 to Kutcher which is similar to the previously described closure but wherein the locking band has several spaced vertical weakening lines formed therein so that upon rotation of the closure during initial removal thereof, the weakening lines rupture dividing the locking band into outwardly flared sectors which remain with the closure upon removal.\nPlastic closures are available which have locking bands which remain connected to the container upon removal of the closure. For example, such a closure is available from Alcoa Corp. under the trademark Drop-Lok.\nShrink seals have the disadvantage of requiring additional assembly operations including the application of heat which is relatively expensive. Moreover, shrink seals are not considered desirable from the point of view of the purchaser since they usually require a separate removal operation to be performed. Closures which leave a locking band on the container are not entirely satisfactory, especially where the containers are returnable and refillable or recycleable, since such bands must first be removed by clippers, knives or the like which is a costly and tedious operation."}
{"text": "1. Field of the Invention\nEmbodiments of the present invention generally relate to a display device and a method of manufacturing the same, and more particular, to a display device including a flexible circuit board with improved reliability and a method of manufacturing the same.\n2. Description of the Related Art\nCurrently, according to advances in display technology, display devices are widely used in portable display devices, such as notebook computers, mobile phones or portable media players (PMPs) as well as home display devices such as TV sets and monitors. With the trend toward the lighter and slimmer display devices, growing attention is focused on a liquid crystal display (LCD) device, an organic electroluminescent display (OELD) device, and so on.\nIn various display devices, including the LCD or the OELD, a flexible circuit board is used as means to be electrically connected to a display panel.\nThe flexible circuit board is easily bendable, compared to a rigid printed circuit board. In the flexible circuit board, a wire pattern is formed on one or both surfaces of a flexibly bendable base film, which can be advantageously used in a space where a general inflexible printed circuit board cannot be installed.\nThe flexible circuit board is generally used in a bent state to provide an electrical connection between a display panel and a circuit board to apply a particular signal to the display panel. Here, since no separate member for supporting the bent portion of the flexible circuit board exists in the flexible circuit board, the flexible circuit board may be deformed. In addition, damages may be caused to the flexible circuit board due to friction between the flexible circuit board and the display panel."}
{"text": "Typical coaxial patchcord assemblies and plug boards are exemplified in U.S. Pat. Nos. 3,284,755, 3,341,801, and 3,335,388. These systems are designed for making board to board interconnections with a front board being preprogrammed with coaxial patchcords point-to-point. The front board is then plugged into a rear board with all of the preprogrammed patchcords being simultaneously interconnected with the boards. The coaxial patchcords are individually provided with very carefully machined connector elements. A typical connnector element, shown in U.S. Pat. No. 3,284,755, includes a smooth noise pin portion connected to the coaxial cable center conductor. The pin is encircled by a sleeve of dielectric material, which, in turn, is encircled by a conducting sleeve, also smooth and having resilient spring properties to allow retention in a corresponding aperture of a front board. Such a connector element is expensive to manufacture and assemble onto a coaxial cable. There has been a long existing need for an economical coaxial connector for point-to-point use without having to resort to patchcord connector elements."}
{"text": "The present invention relates to a slide bearing system surrounding a shaft in a roll for a paper/board machine or a finishing machine, comprising a body section which is formed with a bearing housing, said bearing housing being provided with hydrostatic bearing elements disposed around the shaft and fitted with slide elements, having a sliding surface directed towards the shaft, a bushing between the slide elements and the shaft, which is mounted on the shaft for rotation therewith, said slide elements being adapted to position themselves around the bushing with the sliding surfaces against the external surface of the bushing directed away from the shaft for supporting said shaft, and hence the roll, rotatably relative to the body section, and end caps, along with packing elements therefor, mounted on the axial ends of the bearing housing for sealing the bearing housing. The present invention relates further to a method for slide bearing a roll in a paper/board machine or a finishing machine, by means of an extra-roll slide bearing system.\nThis type of external slide bearing system or assembly can be used in connection with a variety of rolls, such as thermol rolls, center rolls, sizer rolls, felt rolls, etc. FI Patent application 970624 discloses an extra-roll slide bearing arrangement, wherein the slide elements of hydrostatic bearing elements mounted around a shaft are adapted to position themselves around a neck of the roll shaft for supporting the shaft, and hence the roll, rotatably relative to a bearing block. This type of prior art arrangement is depicted schematically in FIG. 1. The prior art solution shown in FIG. 1 comprises bearing elements, which are mounted on a bearing block 5 and capable of positioning themselves freely around a shaft 2 in a nip roll 1. Said nip is established together with a counter roll 1A. Reference numeral 3A represents a main load-carrying bearing and reference numeral 3 represents an abutment bearing. Reference numeral 3B indicates side bearings. Reference numerals 4A, 4 and 4B represent slide elements for bearing elements, respectively. As a result of the conicity of the shaft 2, the bearing elements are positioned to be slightly inclined relative to the rotation axis of the shaft 2, which has required certain special arrangements regarding the bearing block 5. The present Applicant\"\"s earlier Finnish patent 104343 discloses a solution, in which a bushing is fitted between a slide element 4 and a shaft 2, having an external surface substantially cylindrical in cross-section. By virtue of this, the bearing elements can be disposed around the shaft 2 perpendicularly relative to its rotation axis.\nThe assembly and installation of such prior art slide bearings is demanding and involves a multitude of components in the process of fitting a roll with bearings. In addition, these prior art slide bearing systems must be assembled and tested with a roll or a separate test bushing. Moreover, the replacement of such a bearing assembly is a tedious process, and typically the entire roll is replaced by another roll outfitted with a standby bearing assembly. Furthermore, in solutions, wherein sliding surfaces are in a direct contact with a roll shaft, a sliding surface damage means that the roll will be damaged as well.\nConsequently, it is one essential object of the present invention to provide an improved extra-roll slide bearing system, whereby the above drawbacks can be eliminated.\nIn order to achieve this object, a slide bearing system of the invention is characterized in that the slide bearing system is assembled for a single package, which is mountable on a shaft and dismountable from a shaft as a single package. A method according to a first aspect of the invention for slide bearing a roll in a paper/board machine or a finishing machine by means of an extra-roll slide bearing system is characterized in that the slide bearing system is assembled for a single package directly on a shaft. On the other hand, a method according to a second aspect of the invention is characterized in that the slide bearing system is assembled for a single package and mounted as a single package on a shaft.\nOne of the benefits offered by the present invention is that the entire slide bearing system is assembled for a single package, which is readily and quickly mountable on a shaft and dismountable from a shaft and which includes fewer components than the currently available bearing system. In addition, the assembly and testing of a bearing system can be implemented without a roll body and the roll can be replaced without a standby or substitute bearing system. Furthermore, the bushing included in the bearing system can be used as a fitting sleeve for old roll bodies. Still another advantage gained by the invention is an improved repairability in the events of bearing damage."}
{"text": "1. Field of the Invention\nThe present invention relates to a Bi-CMOS logic circuit, and particularly to an integral Bi-CMOS logic circuit in which the drains of an N-type MOS transistor and a P-type MOS transistor are formed as a dual layer diffusing structure to serve as the bases of PNP and NPN bipolar transistors, thereby integrally constructing the functional section and the driving section.\n2. Description of the Art\nAttempts to form different semiconductor devices such as bipolar transistors and CMOS transistors on the same semiconductor substrate have been made for a long time, and a Bi-CMOS semiconductor device has advantages due to a complementing effect such as the high speed and high power driving ability of the bipolar transistor, and the high density and the low power consumption of the CMOS transistor.\nFIG. 1 is a circuit diagram of a Bi-CMOS transistor. In this drawing, the Bi-CMOS logic circuit is an inverter circuit for inverting signals supplied through its input terminal before outputting through its output terminal.\nThe Bi-CMOS logic circuit of FIG. 1 includes: a pair of first transistors, i.e. a P-type MOS transistor MP11 and an N type MOS transistor MN11, with their sources connected respectively to a power voltage Vcc and ground GND, with their gates receiving input signals IN, and with their common drain outputting first output signals OUT11;\na pair of second transistors, i.e. N-type MOS transistors MN13, MN12, with their gate terminals and the base of a bipolar transistor BN11 respectively receiving the input signals IN and first output signals OUT11 of the pair of the first transistors MP11, MN11 respectively, with the drain of one terminal and the source terminal connected respectively to an output terminal OUT and ground, and with the common source terminal and the drain terminal outputting second output signals OUT12; and\na pair of third transistors, i.e. NPN-type bipolar transistors BN11, BN12, with their bases receiving the first output signals OUT11 and the second output signals OUT12 of the second transistors including the N-type MOS transistors MN11, MN12, with their collectors connected respectively to the power voltage Vcc and ground, and with their common emitter outputting the final output signals OUT.\nIn the Bi-CMOS inverter circuit described above, if a high signal is supplied to input terminal IN, the N-type MOS transistor of the pair of the first transistors turns on, and the first output signal OUT11 shifts to a low state. At the same time, the N-type MOS transistor MN13 of the pair of the second transistors turns off, and the MOS transistor MN12 turns on, so that the output terminal OUT is pulled low.\nIf a low signal is supplied to the input terminal IN, the P-type MOS transistor MP11 of the pair of the first transistors turns on, and the first output signal OUT becomes high. When the N-type MOS transistor MN13 of the second transistors turns on, the second output signal OUT12 becomes low. Therefore, the NPN type bipolar transistor BN11 turns on, so that the final output signal OUT is high.\nBi-CMOS logic circuits have been applied to digital logic devices having a high power driving ability, and also applied to the static RAM (SRAM) and to gate array devices.\nThe Bi-CMOS logic circuit described above consists of four MOS transistors and two NPN-type bipolar transistors, with the MOS transistors serving as a functional section, while the bipolar transistors serve as a driving section. Thus, the Bi-CMOS logic circuit above described performs a two-step operation, i.e. inverting the input signals IN and transferring them to the output terminal OUT.\nFIG. 2 is a vertical sectional view of the integrated circuit in which the Bi-CMOS logic circuit of FIG. 1 is formed on an epitaxial layer. That is, an NPN-type bipolar transistor is provided with a buried layer in such a form that: an n- epitaxial layer 22 which functions as a collector region is formed on buried layer 21, and a P base region 24 and an n+ emitter region 23 are formed on the n- epitaxial layer 22.\nAn N-type MOS transistor is formed with n+ source and drain regions 26, 27 on a P- epitaxial layer 25. A P-type MOS transistor is formed with an n- epitaxial layer 29 on an n+ buried layer 28, and P+ source and drain regions 30, 31 are formed on an n- epitaxial layer 29. Here, the n+ regions 32, 34 and the P+ region 33 are the regions for separating the devices.\nReferring to the vertical sectional view of FIG. 2, an exemplary Bi-CMOS logic circuit comprises a functional section consisting of MOS transistors and a driving section consisting of bipolar transistors as described above. Further, they are formed by simply connecting N type and P type MOS transistors and NPN type bipolar transistors which are independent of one another.\nTherefore, there are disadvantages in the above-described structure. First, in terms of the chip, an area as great as the area of the driving section including the bipolar transistors is required, and therefore the entire area of the chip may be increased. Further, since there is a two-stage signal path, the processing speed is slower than that of the inverter circuit which includes the pair of the first MOS transistors MP11, MN11.\nFurther, mutually independent N-type and P-type MOS transistors and NPN-type bipolar transistors are formed on an epitaxial layer of a semiconductor substrate through a large number of process steps and therefore, the process is complicated, and the number of the masks required becomes large, with the result that productivity is reduced.\nIn addition, in the case where the above Bi-CMOS logic circuit is used for obtaining signal products, the formation of vertical PNP transistors becomes impossible, and its application to the actual products is limited."}
{"text": "In general, karaoke television (KTV) is equipped with two display devices, but the two display devices show the same pictures, that is only an MV and lyrics prompts."}
{"text": "Prior art image forming devices include copiers, facsimile machines and printers. A plurality of these image forming devices is connected via network. In the network system with the above connected image forming devices, collaboration may be performed among the image forming devices. For example, the collaboration reduces the operation time by dividing the tasks. For example, when a single original is to be copied one hundred times, the originating image forming device that is given the above task communicates with other image forming devices. As a result of the communication, the task is divided between two image forming devices so that each device copies fifty times.\nIn addition, prior art image forming devices scan the originals on the contact glass and store the image information in memory. Based upon the stored image information, an image is printed on the image transfer paper, and the repeated print operations provide a predetermined number of copies. In the above prior art image forming devices, a remaining memory amount is monitored while the originals are being scanned. The scanning operation is interrupted before the remaining memory reaches zero. The copy operation is determined by the user upon detecting the zero memory. The user can select either terminating the copy operation or starting the copy of the image information that has been already scanned. For example, if an automatic feeding mode is used to feed an original to the contact glass surface from the original tray of the automatic document feeder (ADF) and the remaining memory has reached zero, since some of the originals have been already outputted, the user has to put back the originals in the feeding tray after the user cancels the copy operation and retries the same copy operation.\nFor the prior art technology, the following reference Japanese Patent Publication 2000-69259 discloses that the image information is scanned by an arbitrary image forming device among a plurality of image forming devices, and the image information is transferred to another one of the image forming devices to distribute printing for improved operation in the network system.\nIn the prior art image forming devices, outputted transfer paper is wasted at a child device when a copy job is terminated in response to a memory full detection at the child device. Further, a single original bundle is unfortunately divided for output at a child device when a specified number of copies is outputted based upon the transferred image at the child image. It remains desirable to provide a highly efficient image forming device network system that causes no wasteful miscopy even when memory becomes full at a child device.\nIn the prior art image forming devices, a job is reserved on another machine or a child device before a print job is divided for the transferred image data. The next job reservation is made while the current print job is in operation, is being transferred, or is being printed as a divided job. As Japanese Patent Publication Hei 10-322533 discloses that copiers print out images in the storage unit even if a collaboration unit fails to make printers collaborate during the print operation in the image forming device network. The above printing under abnormal conditions continues as much as possible to prevent decreased productivity by terminating the print operation at other copiers.\nThe Print Operation\nFor the above prior art system, it is necessary to keep the image data in the storage unit so that the print job is continued until the print job is completed at any one of the printers in the system. However, since a child device may have another job in addition to the current collaboration job, to use the memory resources in an efficient manner, it is also necessary to remove image data that no longer requires the print division. For this reason, it remains desirable to provide an image forming device to efficiently utilize the memory resources in sharing a print job among image forming devices that are connected in common network.\nOther prior art references such as Japanese Patent Publications Hei 10-74163 and 2001-13827 disclose the memory resource utilization techniques. To effectively utilize the common memory resources, a memory controller arbitrates the use of the common memory resource among a plurality of software application programs in multifunction devices as well as copiers, printers and fax machines as disclosed in Hei 10-74163. To guarantee a memory area for a collaboration job, a memory block is reserved. However, when there is no collaboration, the memory block is available for any other tasks by a single operation as disclosed in 2001-13827. Furthermore, the rest of the memory block is released if not used even during collaboration. However, the prior art techniques cause restrictions on the memory resources by other application program operations at a child device if the image data is transferred at the beginning of collaboration. For example, if the image data is transferred for collaboration during the execution of a previous job at a child device, the previous job may restrict the memory for a subsequent job. In a worst case, the previous and subsequent jobs compete for memory, and a deadlock may occur. Furthermore, when a next non-collaboration job is to be registered at a child device, the memory resource similarly may not be available. For the above reasons, it remains desirable to control the timing of transferring image data in order to increase the efficiency in the memory resource utilization at a child device.\nThere are other issues that are associated with image data transfers. As described above, Japanese Patent Publication 2000-69259 discloses a prior art technology that the image information is scanned by an arbitrary image forming device among a plurality of image forming devices, and the image information is transferred to another one of the image forming devices to distribute printing for improved operation in the network system. The image data is scanned by an arbitrary master or parent device among the image forming devices. Alternatively, the image data is already scanned and stored in an external memory storage unit. However, when the image data is distributed among child devices, the remaining memory amount may differ among the child devices. The parent memory may reach zero while the original is being scanned as the remaining memory is monitored. As a result, the original is returned, and it is inefficient. For this reason, it remains desirable to increase the productivity in scanning to distribute the image data among the image forming devices for sharing printing. In prior art technologies, remaining memory is monitored while original documents are scanned. Before the remaining memory reaches zero, scanning is interrupted. Then, a user is given a choice to select. He or she can terminate the copying operation. Alternatively, he or she can specify that the printing is initiated for the image data that has been already scanned. In prior art technologies, remaining supplies such as paper and toner are monitored. At the beginning of a job or during the job, a message is generated upon detecting any empty supply for an operator to provide exhausted supplies. The job is continued after supplies are provided. In the image forming device network system, a job is generated by a parent device and shared among child devices in collaboration as disclosed by Japanese Patent Publication 2001-238035. It remains desirable to monitor the supply level at each of the child devices during the assignment of the print job during collaboration.\nIn prior art technologies, a final image including from a parent device and a child device is printed at a child device in order to make the print output to be equal to that from a single image forming device. To provide the final image, it is desired that a new consolidated image is generated at a parent image in stead of transferring the scanned image data to the child image. For this reason, it is not necessary to store the image data at a child device, and the child device memory resource is efficiently utilized. Furthermore, it is desired to transfer the image data to a child device from the last page in the reverse order in case of the stack mode printing at a child device."}
{"text": "The present invention relates, generally, to the field of computer networks, and more specifically, to dispersed or cloud storage of object data.\nObject storage is a computer data storage architecture that manages data as objects, as opposed to other storage architectures such as file systems and block storage. In an object storage architecture, each object typically includes the data, a variable amount of metadata, and a unique identifier. Object storage may be implemented at the device level, the system level, and the interface level. Object storage seeks to enable capabilities not addressed by other storage architectures, for example, application programming interfaces (APIs) that may be directly programmable by an application manipulating data, a namespace that may span multiple physical devices, and data management functions such as data replication and data distribution."}
{"text": "1. Field of the Invention\nThe present invention relates to a liquid crystal display device and, in particular, to a potential generation circuit in an active matrix liquid crystal display device.\n2. Description of the Background Art\nDrive schemes for active matrix liquid crystal display devices are roughly divided into two groups from the viewpoint of a common potential. The first one is generally called a line common inversion scheme where the polarity of the common potential is inverted every other gate line of a liquid crystal display device, and the polarities of the potentials of all data lines are inverted according to the polarity of the common potential. Here, the polarity inversion period may be every other multiple number of lines or may be every frame. The other one is generally called a dot inversion scheme where the common potential is kept constant at a DC potential and the polarities of the potentials of data lines are inverted with reference to the common potential. Although polarity is typically inverted so that adjacent pixels are provided with different polarities, it may be inverted every other multiple number of pixels.\nFor instance, for the case in the line common inversion scheme where the polarity of an applied voltage is inverted every other gate line of a normally white liquid crystal display device, the applied voltage attains a maximum value when the display provides a full black screen, and at that time a current consumed by a potential generation circuit (hereinafter referred to as a “common potential generation circuit”) is also at its maximum. On the other hand, in the dot inversion scheme, when the display provides a full black or white screen, potentials between adjacent pixels are cancelled out, so that a common potential generation circuit consumes little current. However, for example when adjacent pixels are displayed as black and white (typical adjacent pixels have different RGB colors, so adjacent pixels are actually displayed as purple, green, purple, and so on), the common potential generation circuit consumes more power as in the case of the line common inversion scheme. In the dot inversion scheme, polarity is averaged during display of a typical image, so that power consumption by the common potential generation circuit tends to be generally smaller than in the case of the line common inversion scheme.\nConventionally, there is a method that uses a charge recovery circuit in the line common inversion scheme so as to reduce power consumption by the common potential generation circuit that would consume a large amount of power (see, for example, Japanese Patent Application Laid-open No. 2004-69848.) There is also another method for reducing power consumption by the common potential generation circuit (see, for example, Japanese Patent Application Laid-open Nos. 2007-093696, 11-194320, 2001-282197, and 06-088955.)\nSuch conventional common potential generation circuits either fail to achieve sufficiently low power consumption or achieve low power consumption but only with a complicated circuit configuration."}
{"text": "Among the light stabilizers enjoying great commercial success is bis(2,2,6,6-tetramethylpiperidin-4-yl) sebacate, TINUVIN 770, available from CIBA-GEIGY Corporation. The preparation of this sebacate ester is described in U.S. Pat. No. 3,640,928 by the reaction of 2,2,6,6-tetramethylpiperidin-4-ol with an appropriate sebacic acid derivative.\nIsolation and purification of the desired sebacate ester by conventional crystallization methods leads to a product obtained as friable thin flakes with severe dusting problems. While this dusting propensity in no way affects adversely the light stabilization efficacy of this sebacate ester, it does lead to severe handling problems, potential health hazards and related economic barriers to the full realization of the economic potential of this effective light stabilizer in practical use."}
{"text": "The value of many cloud services lies in the ability to operate on client's data. Typical examples include data storage, webmail service, advertising, geolocation services, etc. However, the need to access client's data, e.g., files, email, location is often at odds with clients' data privacy needs, which require encryption to protect their data. Reconciling these contradictory requirements, and achieving computation on encrypted data, is an important research and engineering problem.\nWe focus here on the specific case of indexing and searching on documents. Typically to perform a search on a set of documents, one creates an index of keywords included in the documents and later on performs the search on the index. This does not raise major issues if the querier, who is defined as the one who requests the search, and the storage provider, who is defined as the one who stores the documents and the index and who actually performs the search and returns the result, fully trust each other. The difficulties arise when the querier wants to protect the privacy of its data and store the documents as well as the indexes encrypted, and send encrypted queries as well. In that case, it would be advantageous to be able to perform encrypted searches on a collection of encrypted documents. And more precisely, it would be advantageous to be able to perform all kind of encrypted searches on a collection of encrypted documents, including any Boolean search, such as conjunctive, disjunctive and negative searches.\nBy “disjunctive search”, we mean a search allowing the user to search for encrypted documents containing word1 or word2 . . . or wordn. By “conjunctive search”, we mean a search allowing the user to search for the encrypted documents containing: “word1 and word2 . . . and wordn”. By “negative search”, we mean a search allowing the user to search for all encrypted documents which do not contain one or several particular words.\nCurrently, there is no solution enabling to perform Boolean search on encrypted data, i.e. conjunctive, disjunctive and negation search option at the same time."}
{"text": "1. Field of the Invention\nThe present invention relates to a drive circuit and a power supply apparatus having a switch element such as MIS-FET (Metal Insulation Semiconductor—Field Effect Transistor), IGBT (Insulated Gate Bipolar Transistor) or the like, more particularly to an improved technique for switching the switch element at a high speed.\n2. Description of the Related Art\nGenerally, in drive circuits having a voltage-driven type switch element such as a MIS-FET, an IGBT or the like, controllability of the drive circuits can be improved by driving a gate voltage at a high speed, thereby shortening a switching time of the switch element.\nFurther, when the drive circuits receive a negative voltage of a pulse signal, a method is performed which applies a reverse bias voltage to the switch element in order to abruptly turn off the switch element.\nFrom the above view point, a variety of techniques have been disclosed relating to drive circuits for a switch element.\nFor example, Japanese Unexamined Patent Publication (Kokai) No. H08(1996)-149796 discloses a technology of drive circuits having a capacitor in series between a secondary winding of a transformer and a switch element, thereby applying a reverse bias voltage generated by the capacitor to a gate of the switch element, and limiting amplitude of the reverse bias voltage by a Zenor diode connected in parallel with the capacitor.\nJapanese Unexamined Patent Publication (Kokai) No. H11(1999)-145810 discloses a technology of drive circuits connecting a capacitor and a diode in series between a secondary winding of a transformer and a voltage-driven type switch element, and connecting a base and an emitter of a discharging transistor to both terminals of the diode without a current-limiting resistor, thereby turning on the discharging transistor at a high speed when the switch element is off-state, which allows the gate capacitor to be discharged at a high speed.\nAccording to the technology disclosed in Japanese Unexamined Patent Publication (Kokai) No. H08(1996)-149796, a reverse bias voltage is generated for applying the gate of the switch element. However, when an input voltage is reversed and a gate capacitor of the switch element is discharged, a large current flows through the gate, the capacitor and the transformer (inductance). Therefore, it suffers from a disadvantage that a large level of a LC resonance appears.\nMoreover, when an input voltage is reversed, a time constant of the discharge is long in a discharging loop formed by the gate (gate capacitor), the capacitor and the transformer (inductance). Therefore, it also suffers from a disadvantage that power dissipation (consumption) is large.\nAccording to the technology disclosed in Japanese Unexamined Patent Publication No. H11-145810, a current-limiting resistor is not provided and impedance of a discharging loop from the gate capacitor is small when the switch element is off-state. As a result, it suffers from a disadvantage that a large level of a LC resonance may appear and power dissipation is large because of a large current that appears due to the small impedance.\nFurther, since a power loss of the discharging transistor becomes large under high temperature range, power dissipation may be even larger."}
{"text": "In the context of radiation therapy a target within a human body generally has to be irradiated in order to treat a disease for example, in particular a tumor. In this process a radiation therapy apparatus applies a high radiation dose in a targeted manner in an irradiation region (isocenter) of the human body. During irradiation the problem frequency occurs that the irradiation region moves and/or is displaced. Thus a tumor in a stomach region for example moves during an respiratory process of a patient. On the other hand a tumor can grow or reduce in size in a time interval between radiation planning and the actual irradiation.\nIt has therefore been proposed to monitor a position of an irradiation target in a patient during irradiation by means of medical imaging. This allows a beam and/or beam guide for irradiation to be controlled or optionally irradiation to be terminated. It is also of significant interest to track a radiation focus in respect of the actual position of the irradiation region.\nA combined radiation therapy apparatus and magnetic resonance apparatus is particularly advantageous. This features high soft part resolution compared with a computed tomography apparatus for example, so that an advantageous contrast can be shown in this region.\nFor efficient irradiation a radiation source of the radiation therapy apparatus is positioned as close as possible to a patient. To this end the radiation source is generally disposed at least partially within the magnetic resonance apparatus and in particular within a magnetic field of the magnetic resonance apparatus. However this arrangement has the disadvantage that an electron path of electrons of an electron beam of the radiation therapy apparatus is subject to interference from the magnetic field of the magnetic resonance apparatus.\nWO 03/008986 A2 proposes a separation of gradient coils of the magnetic resonance apparatus and a tailored design of a main magnet, so that an almost magnetic-field-free space is produced outside the magnetic resonance apparatus. However this arrangement has the disadvantage that the apparatus is of large extension and only one angle of incidence is available for a radiation treatment. Also the separated gradient coils result in significant disadvantages in respect of the image quality of medical imaging. Also the radiation from the radiation therapy apparatus must penetrate a steel body of the magnet, resulting in a deterioration and/or degradation of the beam profile and beam intensity.\nA combined magnetic resonance apparatus and radiation therapy apparatus, in which however an x-ray and/or gamma beam is generated for the purposes of irradiation outside the magnetic resonance apparatus and therefore outside the active region of a magnetic field, is also known from U.S. Pat. No. 6,198,957 B1. This means that the x-ray beam is generated a very long distance from the actual treatment region, so the apparatus is likewise of large extension, in particular with a variation of an angle of incidence. Also the long distance means that a high radiation dose has to be generated, to achieve the required penetration depth of radiation for the radiation treatment.\nA guide for an electron beam along a main axis of the magnetic resonance apparatus is also known from DE 102008007245 A1. The electron beam is deflected through 90° to collide with a target. An electron beam and the target are disposed within a patient support of the magnetic resonance apparatus. However this means that the space available for the patient within the magnetic resonance apparatus is additionally limited by the radiation therapy apparatus."}
{"text": "1. Field of the Invention\nThis invention relates to the removal and recovery of water soluble organics (WSO) from oil process water. In particular, the present invention relates to the removal and recovery of certain water soluble petroleum organics from crude oil production water and from aqueous streams used in removing water and/or inorganic salts, such as sodium chloride, from crude oil, residual oil, waste oils and the like.\nIn the removal of crude oil or production fluid from earth formations significant quantities of water are quite often removed along with the oil. In the Middle East, the production fluid can be virtually pure oil. However, it is not uncommon that oil well production fluids are composed of 90% or more of water and only 10% or less of crude oil. Such water is referred to as connate water and is produced along with the oil. Oil well production fluid typically comprises water, oil, water soluble substances (e.g., naphthenic acids, phenols and inorganic salts) and cations that tend to form scale (e.g., calcium, barium and strontium ions).\nOne of the first steps after removal of the oil well production fluid is to separate the oil from the water by phase separation techniques. The phrase \"oil process water\" refers to the water of such oil water production fluid as well as to the water phase produced by such phase separation.\nConventionally, separation is accomplished using a bulk separator or a free water knock out system. Virtually all of the hydrocarbon is conveniently recovered in this manner. Unfortunately, certain organic compounds, as well as inorganic salts and acids, are soluble in water; and mere phase separation will not remove the water soluble compounds from the aqueous phase. Water soluble organics include, among other things, certain naphthenates, phenols, lower fatty acids, etc. Water soluble inorganic salts include sodium chloride, sodium sulfate, calcium chloride, barium chloride, etc. While the amount of water soluble organics may be relatively small, up to 1,000 ppm, they nevertheless give rise to environmental problems, when the aqueous phase is discharged into the environment without the removal of the water soluble organics. Furthermore, the water soluble organics may be valuable substances. In order to meet present day strict environmental standards, a process to reduce the level of the water soluble organics in the discharged streams to 50 ppm or less is needed. The present invention provides a simple, economical procedure for accomplishing this end. While known methods can reduce the content of water soluble organics to the desired low level, they are relatively complex and/or expensive in comparison with the process of the present invention. For example, one could separate the water and water soluble organics by distillation or the use of biological treating ponds.\nUse of certain petroleum fuel oils for specific purposes requires that such oils be treated in order to remove undesirable corrosive contaminants therefrom. For example, fuel oils used in the newer, high efficient, gas turbine power plants must meet certain strict requirements with respect to the presence of inorganic contaminants in the oil. The presence of sodium chloride and other similar inorganic salts renders the oil less than suitable as a fuel for use in gas turbine power plants. To upgrade fuel oil so that it is acceptable for use in the turbines, the fuel oil is commonly processed using a multi-stage electrostatic desalting facility. In such operation essentially complete removal of water soluble inorganic salts from fuel oil is accomplished by a counter-current water washing process in which a high electrostatic gradient is used to break the interim water-in-oil emulsion. Selective specialized chemical demulsifiers are normally used in such salt removal processes. While the process removes the inorganic salts from the fuel oil, unfortunately the wash water removes the water soluble organics from the fuel oil. In view of strict environmental standards, it is of considerable importance to remove these water soluble organics present in the wash water of the inorganic salt removal processes used to render the fuel acceptable for burning in turbines.\nThe invention is directed to a simple, straightforward method of removal of water soluble organics from oil well production fluids, as well as from aqueous streams used to render fuel oil acceptable for use as fuel for gas turbine power plants by the removal of inorganic salts therefrom.\n2. Description of the Prior Art\nU.S. Reissue Pat. No. 29,908 discloses recovery of oily waste in water by the use of a demulsification-flocculation separation process employing a combination acid-alum treatment. This process has no effect of removing water soluble oils.\nU.S. Pat. No. 3,687,845 discloses removal of tramp oils from oil-in-water emulsions by the addition of a water soluble anionic polymer of high molecular weight and heat treatment of the emulsions so that the oils are coalesced and float to the surface of the liquid. Water soluble oils are unaffected by this treatment.\nU.S. Pat. No. 3,707,464 discloses the addition of caustic or acid to adjust the pH of the liquid to about 8.0 to accelerate gravitational separation of oils and solids and the use of elevated temperatures to accelerate the separation process. Rather than being removed, water soluble organics in the solution would be stabilized by this treatment.\nU.S. Pat. No. 4,035,289 discloses the use of microorganisms to reduce the presence of organic material in mineral oil drilling fluids. This is an example of the slow, expensive know methods involving biological treatments or use of absorbent materials, like charcoal and/or ion exchange resins.\nU.S. Pat. Nos. 4,818,410 and 4,839,054 are directed to techniques for removing water soluble organics from oil process water by acidifying the water, contacting the acidified water with oil and thereafter separating the oil from the water. However, these patents recommend use of a mineral acid such as sulfuric acid, hydrochloric acid or phosphoric acid. Use of such acids involves in several drawbacks. Sulfuric acid and hydrochloric acid are corrosive to metal surfaces in the water to be treated. Phosphoric acid forms a precipitate with multivalent cations such as Ca.sup.+2 in the water, resulting in formation of scale and oil-wetting particles that tend to retain oil in the water phase. The agglomeration of oil onto such particles reduces the quality of phase separation and the particles carry the oil through the system, thereby fouling the system."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor device, and more particularly to a semiconductor device in which a front-end unit having an interface function and a back-end unit including a memory core are integrated on separate semiconductor chips.\n2. Description of Related Art\nSemiconductor devices such as a DRAM (Dynamic Random Access Memory) are often used in a state where a plurality of the semiconductor devices are mounted on a module substrate. The modularized DRAMs are sometimes classified into a plurality of ranks that are exclusively selected by chip select signals (see Japanese Patent Application Laid-open No. 2010-134904). Because different ranks can be accessed independently of each other unless competition on a data bus does not occur, use efficiency of the data bus can be enhanced by classifying the DRAMs on the module into the plural ranks.\nOn the other hand, there has recently been proposed a technique in which a so-called front-end unit performing an interface with a memory controller and a back-end unit including a memory core are integrated on separate chips, and these chips are stacked to form a single semiconductor memory device (see Japanese Patent Application Laid-Open No. 2007-158237). According to this technique, in a core chip on which the back-end unit is integrated, a space that can be allocated to the memory core increases, whereby a storage capacity per 1 chip (per 1 core chip) can be increased. On the other hand, the interface chip on which the front-end unit is integrated can be manufactured by a process different from the process for the memory core, whereby a circuit can be formed with a high-speed transistor. Furthermore, plural core chips can be allocated to one interface chip, resulting in that a high-speed semiconductor memory device having extremely large capacity as a whole can be provided.\nHowever, in the stacked semiconductor device, through silicon vias used to transmit commands are commonly connected between the plural core chips. Accordingly, when the plural stacked core chips are classified into plural ranks, a transfer interval of the commands on the through silicon vias is reduced and thus the commands may not be correctly transferred. To prevent this, an issuance interval of the commands needs to be extended to avoid competition of the commands also at the time of accesses to different ranks; however, use efficiency of the bus is lowered in this case."}
{"text": "The present disclosure relates to a sheet conveying device for conveying a sheet member, and an image reading device and an image forming apparatus including the sheet conveying device including a conveying roller configured to be rotatable in a forward rotational direction and a reverse rotational direction.\nThe image reading device such as a scanner is provided with an automatic document feeder (hereinafter referred to as “ADF”) for automatically conveying a document sheet set in a sheet feeding tray. The ADF includes a feeding roller for feeding inside the document sheet set in the sheet feeding tray, and a discharging roller for discharging the document sheet to an external discharge tray. The document sheet conveyed to a reading position by the feeding roller is discharged to the discharge tray by the discharging roller.\nIn an image reading device including the ADF as described above, conventionally, there is known a stopper member for preventing the document sheet in the discharge tray from being drawn into the discharging roller, when the discharging roller is rotated in a rotational direction opposite to that at a time of discharge. The stopper member is configured to prevent a rear end of the document sheet from coming in contact with a surface of the discharging roller for discharging the document sheet."}
{"text": "Cyclin-dependent kinases (CDKs), belong to a group of serine/threonine kinases involved in the regulation of cell cycle progression, neuronal function, differentiation and apoptosis. Their activity is tightly regulated by multiple mechanisms including binding to corresponding cyclins, of which level of expression oscillates throughout the different phases of the cell cycle. Different CDK/cyclin complexes are activated during each cell cycle step through G1, S, G2, M phases. Sequential phosphorylation of the retinoblastoma protein (pRb) by CDK4/cyclin D, CDK6/cyclin D in early G1 phase and CDK2/cyclin E in late G1 phase causes the release of the E2F, proteins of transcription factor. In turn, E2F proteins lead to transcriptional activation of a set of genes required for entry into S-phase of the cell cycle. CDK2 is subsequently activated by cyclin A during the late stages of DNA replication, S-phase, and promotes appropriately timed deactivation of E2F to prevent apoptosis triggered by persistent E2F activity. Finally, CDK1 in complex with cyclin A or B is thought to be have roles in regulating the G2/M checkpoint and driving cells through mitosis.\nIn addition to the cell cycle control, other roles have been determined for CDK2, 7, 8 and 9. For example, CDK2/Cyclin E is important to the p53 mediated DNA damage response pathway and CDK7, 8 and 9 are involved in the regulation of transcription initiation and elongation through phosphorylation of RNA polymerase. Therefore, CDKs affect cell growth and survival through several different mechanisms and proper regulation of CDK activity is important to various cellular processes. It is now recognized that deregulation of CDKs by abnormal high expression of cyclin such as cyclin D and cyclin E or mutation occurs in many human tumors. For example, the expression and catalytic activity of CDK2/cyclin E complexes is increased in colorectal, ovarian, breast, and prostate cancers and elevated expression of cyclin E in primary tumors has correlation with poor survival rates for breast cancer patients. Abnormal expression of CDK1/cyclin B complexes has been also observed in some cases, prostate adenocarcinoma and lung cancer.\nAlthough a report have shown that CDK2 may not be required for cell cycle progression and proliferation, recent reports suggested that melanocytes and hepatocytes may be dependent on CDK2 for proliferation and survival. Also, an investigation of simultaneous depletion of CDK1 and CDK2 was reported to provide increased efficacy in anti-proliferation of tumor cell lines, compared with targeting either CDK1 or CDK2 alone. In addition, emerging evidences indicate that certain tumor cells may require specific interphase CDKs for proliferation. Thus inhibition of CDKs may provide an effective strategy to control tumor growth as attractive targets for cancer therapy.\nTo date, a number of small-molecule CDK inhibitors are currently under clinical trials. These inhibitors are flat, small heterocycles which act by competition with ATP in the kinase ATP-binding site. Among them, flavopyridol, was the first CDK inhibitor to enter clinical evaluations. R-Roscovitine (trisubstituted purine analog) and BMS-387032 (aminothiazole) are selective for CDK2/cyclin E and PD-0332991 (pyridopyrimidine) is highly selective for CDK4/cyclin D and CDK6/cyclin D. Indirubin, a bis-indole scaffold and its derivatives have been investigated with considerable interests as potent inhibitors targeting important protein kinases such as CDK, GSK-3β, and aurora kinases.\nIn WO 2005/070416 A1, the inventors of present application has disclosed ‘indirubin derivatives having anticancer property against human cancer cell line’. In this disclosure, indirubin derivative having anticancer property by inhibiting cell proliferation as to human cancer cell line has been disclosed. Among disclosed indirubin derivatives, 5′-bromo-5-nitro-indirubin-3′-oxime showed the most potent anti-proliferative effects (IC50=0.79˜2.9 μM) against various cancer cell lines. However, the further substitution of radicals at the 5′ position in indirubin skeleton has not been made as well as anti-cancer activity regarding these compounds has not been measured. Further, there has been no report regarding the synthesis and biological evaluations exploring the effect of various substitutions at the 5′ position of indirubin skeleton.\nOn the other hand, the inventors of present application have found that the additional substitution at the 5′ position is clearly favorable compared with 5-nitro-indirubin-3′-oxime analog. Therefore, the inventors of present application have completed the invention by design, synthesis and biological evaluation including CDK inhibitory activity, anti-proliferative activity and in-vivo anti-cancer activity of novel 5′,5-substituted indirubin-3′-oxime analogs."}
{"text": "The present invention relates to compositions and processes for applying strong protective acrylic coatings to fingernails and toenails with or without artifical extensions to the nails.\nNumerous techniques for applying protective coatings for fingernails, with or without extensions to the nails, exist today. The types of artificial nails produced by these techniques can be classified as (1) glue-on nails, (2) nail wraps, (3) sculptured nails, and (4) nail dips.\nGlue-ons are pre-formed plastic fingernails having the desired, usually extended shape of the natural nail. These are glued directly on to the fingernail. Because of the wide variety of nail shapes, glue-on nails are difficult to fit. As the nails grow, the glue-on nails move distally creating an unsightly ledge across the base of the nail. This is repaired either by removing the artificial nail and attaching a new one or by filling in behind the ledge by a sculpturing technique as described below. The former repair is expensive because it must be performed every two to three weeks and the removed nail cannot be reused. The latter repair requires the skills of a scupltured nail technician.\nIn the nail wrap method a piece of paper or fabric is cut to the shape of the natural nail, glue is applied to the natural nail and the paper or fabric is pressed into the glue. The paper or fabric is then smoothed insofar as possible and folded over the free edge of the nail. The folds are glued to the underside of the nail and more glue is applied to the paper or fabric in successive layers until the surface is smooth. This type of artificial nail has many problems. First of all, the nail is made if glues, usually the lower cyanoacrylate esters, which do not produce a particularly strong nail. The single lamina of paper or cloth adds little strength to the composite. Secondly, the surface of the natural nail is a complex curve and it is difficult for the paper or fabric to conform to the surface without forming folds. Under the nail, the folded over paper or fabric forms convolutions which are unattractive when viewed from the edge or from underneath. Moreover, these folds harbor dirt and microorganisms. If edge folds are not used or if the nails are trimmed, the edge is ragged and readily delaminates. As with glue-ons, the artificial nails move distally as the nail grows and must be filled in or the artificial fingernail must be removed and replaced with a new nail wrap.\nNail sculpting has become popular with those who are unsatisfied with the methods described above. However, nail sculpting requires a great deal of training and skill. In this method, preparation of the natural nail is important because the thick, tough acrylic nail becomes permanently attached to the point where it is applied. The nail is first cleaned, roughened, and treated with a bromide. It is then coated with glue, usually a cyanoacrylate ester, or a primer containing methacrylic acid which binds the acrylic to the natural nail. Then the sculpturist wets a small brush with a liquid consisting of a blend of mono-, di- and trimethacrylate esters and a promoter, usually N, N, dimethyl-para-toluidine which is capable of inducing the decomposition of benzoyl peroxide at room temperature. The wetted brush is then dipped into a finely powdered polymethacrylate ester which contains benzoyl peroxide. The liquid and powder form a dough which begins to cure (i.e., harden or polymerize) immediately. The dough is then quickly placed onto the prepared nail and smoothed or sculptured in place. More dough is added until the nail is covered and filled to the desired, uniform thickness and shape. Filing, buffing and, sometimes, the addition of more dough result in a strong, albeit thick nail ready for polish. If the fingernail is to be extended, a form with a non-adhering upper surface is placed under the edge of the nail and the sculpturing is continued out over the form. The form is removed after the dough is cured.\nExtraordinary skills are required for this procedure. The dough must be of a precise consistency. If it is too thin, it is runny and cures slowly and sometimes does not cure at all. If it is too thick, it is weak due to air bubbles and stress cracks. If the sculpturist pauses too long between applying each brush full of dough, the bond is poor and the nail cracks easily. If the dough is not placed and worked skillfully an excessive amount of filing is required to remove humps and lumps to produce an attractive nail.\nRecently, a method of applying arificial fingernails called nail dipping has been introduced. In this method, the natural nail is coated with a high-viscosity glue of the cyanoacrylate ester type. The wetted nail is then briefly dipped into a powdered, uncross linked polymethacrylate ester. Any excess powder is brushed off and the surface is smoothed partially by filing and light buffing. This rough surface is then sealed or filled in by applying a low-viscosity cyanoacrylate ester glue in one or more coats until the powder no longer protrudes above the surface. The surface is then filed and buffed and is ready for the application of nail polish.\nIt is apparent that this method circumvents the fitting problems of glueing-on and wrapping and the skills required of sculpting. However, the finished nail is not strong. The cyanoacrylate ester glues, although adherent, are not structurally sound. Layers as thick as artificial nails if not supported by opposing substrates break easily. The uncross linked polyethylmethacrylate powder is weak per se and lends little toughness to the polycyanoacrylate ester surrounding it. Thus, the dip process nail may be stronger than nail polish but it is much weaker than a sculpted acrylic nail.\nIt follows, therefore, that there is a present need for strong, easily appliable, repairable, fillable, artificial nails and a better procedure to repair and strengthen natural nails and to apply the artificial nails."}
{"text": "Generally described, computing devices and communication networks can be utilized to exchange information. In a common application, a computing device can request content from another computing device via a communication network. For example, a user at a personal computing device can utilize a software browser application to request a Web page from a server computing device via the Internet. In such embodiments, the user computing device can be referred to as a client computing device and the server computing device can be referred to as a content provider.\nContent providers are generally motivated to provide requested content to client computing devices, often with consideration of efficient transmission of the requested content to the client computing device and/or consideration of a cost associated with the transmission of the content. For larger scale implementations, a content provider may receive content requests from a high volume of client computing devices, which can place a strain on the content provider's computing resources. Additionally, the content requested by the client computing devices may have a number of components, which can further place additional strain on the content provider's computing resources.\nWith reference to an illustrative example, a requested Web page, or original content, may be associated with a number of additional resources, such as images or videos, which are to be displayed with the Web page. In one specific embodiment, the additional resources of the Web page are identified by a number of embedded resource identifiers, such as uniform resource locators (“URLs”). In turn, software on the client computing devices typically processes embedded resource identifiers to generate requests for the content. Often, the embedded resource identifiers reference computing devices associated with the content provider such that the client computing devices would transmit requests for resources to the referenced content provider computing devices.\nSome content providers attempt to facilitate the delivery of requested content, such as Web pages and/or resources identified in Web pages, through the utilization of a content delivery network (“CDN”) service provider. A CDN service provider typically maintains a number of computing devices in a communication network that can maintain content from various content providers. In turn, content providers can instruct, or otherwise suggest to, client computing devices to request some, or all, of the content provider's content from the CDN service provider's computing devices.\nCDN service providers are also generally motivated to provide requested content to client computing devices often with consideration of efficient transmission of the requested content to the client computing device or consideration of a cost associated with the transmission of the content. Typically, CDN service providers often consider factors such as latency of delivery of requested content in order to meet service level agreements or generally improve the quality of delivery service. In turn, the CDN service provider can utilize the considered factors in processing resource requests."}
{"text": "1. Field of the Invention\nThe present invention relates to servo systems and, in particular to a new and improved servo speed control circuit.\n2. Prior Art\nIn the past, servo systems have been designed with linear gain to drive a servo motor at a fixed speed. In particular, such prior art servo systems having a feedback circuit for speed control employ a constant current generator to provide a linear ramp voltage for sampling purposes. The servo system is then tuned for the desired gain and phase margin at a selected speed.\nA servo system tuned for operation at one speed may be unstable at other speeds since the gain required in the feedback circuitry is different for different speeds. Thus, if such prior art circuits are used to drive the servo motor at more than one speed, additional circuitry is required to control the speed of operation at each of the desired speeds.\nAnother approach to the above problem is to employ different current generators adjusted for each of the desired speeds, which generators are selectively switched into the feedback circuit so as to provide the different feedback gains for these different speeds. A problem with this approach is that the gain is correct only at those speeds having the appropriate gain. Other problems are that such circuitry is unreliable, cumbersome and complex."}
{"text": "This invention relates to an adapter for a discharge lamp. The adapter is particularly suited for use with a discharge lamp in which a discharge tube is disposed substantially in a plane and a lamp support housing is substantially surrounded by the discharge tube.\nU.S. Pat. No. 4,241,386 discloses an adapter for a discharge lamp having a circular discharge tube, which surrounds a circular zone. The discharge lamp has a plug portion disposed on the peripheral region of the circular zone. Accordingly, this known adapter also comprises a socket portion matching the plug portion of the circular discharge lamp. The socket portion must be positioned in a peripheral region of the adapter in order to accommodate the circular discharge lamp.\nMost mounting sockets receiving such circular low pressure discharge lamps (also called fluorescent lamps) have a socket portion which is situated in a peripheral area of the mounting socket, corresponding to the peripheral position of the plug portion of the circular discharge lamp.\nU.S. Pat. Nos. 4,549,251 and 4,458,301 disclose a discharge lamp with a discharge tube where the discharge tube is bent in a shape so as to surround a substantially planar zone. The discharge lamp includes a lamp support housing. The lamp support housing is positioned in an substantially central area of the zone defined by the discharge tube. The lamp support housing comprises a plug portion with contact pins providing the electrical connection of the discharge lamp.\nIt would be desirable if discharge lamps of the type described in U.S. Pat. Nos. 4,549,251 and 4,458,301 were connectable to mounting sockets originally designed for circular lamps, i. e. for mounting sockets which have their socket portion in a peripheral area.\nU.S. Pat. No. 5,471,375 teaches a mounting socket for use with a discharge lamp having a similar layout as those described in U.S. Pat. Nos. 4,549,251 and 4,458,301. The mounting socket comprises a socket portion matching the electric contacts of the discharge lamp, and also comprises a plug portion matching a known screw-type socket for incandescent lamps. Both the socket portion and the plug portion of the mounting socket are positioned in a central area of the discharge lamp when the discharge lamp is inserted in the mounting socket. However, when the mounting socket is placed in the screw-type socket of the incandescent lamp, and the discharge lamp is inserted in the mounting socket, the plane of the discharge lamp socket is above the level of the screw-type plug portion of the mounting socket. For this reason, a discharge lamp inserted in the mounting socket may interfere with the proper positioning of a lamp cover originally designed to cover a lamp inserted in the lower socket. Thus there is a particular need for a structure in which the level of the discharge lamp plug portion is not much higher than the plug portion of the mounting socket, and the plane of the discharge lamp is substantially the same as the plane of the original socket.\nIn an exemplary embodiment of the invention, an adapter for a discharge lamp comprises a socket portion suitable for receiving electrical contacts of a first discharge lamp. The first discharge lamp includes a discharge tube, which is disposed substantially in a planar zone having a principal plane. The socket portion has an associated socket contacting plane. The associated socket contacting plane is defined as the principal plane of the discharge tube of the first discharge lamp, when the latter is inserted into the socket portion of the adapter.\nThe adapter also comprises a plug portion which is insertable into an external female socket. The external female socket is of the type which is suitable for receiving a second discharge lamp, where the second discharge lamp has a discharge tube with a planar tube configuration. The plug portion of the adapter has an associated plug contacting plane. This plug contacting plane is defined by the principal plane of the tube configuration of the second discharge lamp, when the latter is inserted into the external female socket. The adapter further comprises a connecting portion for connecting the socket portion and the plug portion. The connecting portion provides a substantially rigid connection between the socket portion of the adapter and the plug portion of the adapter so that the socket contacting plane substantially coincides with the plug contacting plane.\nThe adapter is particularly suitable for discharge lamps which have a discharge tube surrounding a lamp support housing so that the lamp support housing is positioned in an substantially central area of the plane defined by the discharge tube. Such a lamp will be termed, folded lamp,, hereinafter because the ends of the discharge tube are folded inwards into the lamp support housing. The connecting portion of the adapter may be sized so that the plug portion is positioned in an substantially peripheral area of the plane defined by the discharge tube when the plug portion of the discharge lamp is inserted in the socket portion of the adapter. This allows discharge lamps with central plug portions to be used with mounting sockets and lamp housings which were designed for circular discharge lamps. Due to the fact that the plug contacting plane of the adapter coincides with the socket contacting plane, the principal plane of the discharge tube of the folded lamp will be practically in the same plane as the plane of the circular lamp. This means that the lamp covers originally designed for circular discharge lamps will be readily applicable with folded lamps as well."}
{"text": "1. Field of the Invention\nThis invention relates to media containing an indicator, 2,3,5-triphenyltetrazolium chloride, and a selective agent composition; and methods for the recovery and enumeration of Campylobacter species.\n2. Description of the Related Art\nCampylobacter species have been recognized as important causative agents of foodborne illness. There is a strong association of foods of animal origin in the transmission of disease to humans. Poultry is one such food with high carriage rates of Campylobacter contamination. Campylobacter jejuni, C. coli and C. lari are known to cause an estimated 2.2 million cases of foodborne gastroenteritis per year in the United States alone (Tauxe et al., American J. Public Health, Volume 77, 1219-1221, 1987). The vast majority of these cases are associated with the consumption of improperly prepared or handled foods. Although the origin of this disease in humans is primarily linked to poultry, the food microbiology and poultry communities have been slow in directing substantive attention toward the organism. This has been due, in part, to the unique physiological requirements of these organisms, impairing their culture and identification from foods and clinical specimens.\nA variety of enrichment and culture media have been proposed for the isolation of Campylobacter species (Park et al., Campylobacter, In: Compendium of Methods for the Microbiological Examination of Foods, second ed., M. L. Speck (ed.), Am. Pub. Hlth. Assoc., Wash., D.C., 386-404, 1984xe2x80x94the contents of which are herein incorporated by reference). Because Campylobacter can be overgrown by other organisms present in sources, the use of selective media, incorporating antibiotics and/or antimicrobial agents, is essential for their isolation. Ideally, any culture medium selected should also be differential, allowing the characterization of the Campylobacter by distinctive colonial appearances in culture.\nRapid and sensitive methods for recovering Campylobacter would be useful for both epidemiological work and routine examination of food sources. The main drawback associated with numerous available procedures is the length of time needed for enrichment. Enrichment culture incubation ranges from 16 to 48 hours before plating on selective media, which then requires an additional 24 to 48 hours for isolation. This 3 to 4 day procedure is difficult to reconcile with rapid marketing strategies while maintaining interest in the public health.\nThe unique physiological requirements of Campylobacter species provide difficulties in culturing the microorganism. C. jejuni require special microaerobic atmospheres for growth (Kiggin et al, J. Bacteriology, Volume 72, 397-400, 1956), and its translucent colonies are frequently difficult to identify on dark, opaque Campylobacter agars.\nSeveral agar media have gained prominence for the isolation of Campylobacter. Campy-Brucella Agar Plate (Campy-BAP), has been widely used and cited in the Compendium of Methods for Microbiological Examination of Foods (2nd Ed., American Public Health Association, Wash., D.C., M. L. Speck, ed., 386-404, 1984). Campy-BAP agar contains Brucella agar, lysed horse blood, vancomycin, polymyxin B, trimethoprim lactate, amphotericin B and cephalothin. Food samples assayed with Campy-BAP medium often yield large numbers of breakthrough flora. Butler developed a selective medium for C. jejuni containing a nutrient agar base, blood, and five selective agents, cycloheximide, cefazolin, bacitracin, colistin sulfate and novobiocin as described by Smibert (Campylobacter, In: Bergey\"\"s Manual of Systematic Bacteriology, Krieg and Holt (eds.), Williams and Wilkins, Baltimore, Md., Volume 1, 111-115, 1984; the contents of which are herein incorporated by reference). Another agar, Campy-Cefex, allows for selective and differential culture of C. jejuni (Stern et al., Journal of Food Protection, Volume 55 (7), 514-517, Jul. 1992; U.S. Pat. No. 5,891,709, Apr. 6, 1999). Campy-Cefex agar contains Brucella agar, 0.05% ferrous sulfate, 0.02% sodium bisulfite, 0.05% sodium pyruvate, 33 mg/L sodium cefoperazone, 200 mg/L sodium cycloheximide, and 5% lysed horse blood. CCDA agar contains Nutrient broth No. 2, Bacteriological charcoal, casein hydrosylate, sodium desoxycholate, ferrous sulfate, sodium pyruvate, agar, yeast extract, sodium cefoperazone, and sodium cycloheximide. CCDA medium, also widely used and cited in Compendium of Methods for Microbiological Examination of Foods, (supra), was developed to replace the blood component which is specified in many Campylobacter recovery media, with charcoal. It uses cefoperazone as the selective antimicrobial agent acting in concert with a 42xc2x0 C. incubation temperature and microaerobic atmosphere to limit the proliferation of non-Campylobacter organisms. The 42xc2x0 C. incubation temperature greatly reduces the need for anti-gram-positive antimicrobials. The ferrous sulfate component of the medium has been used to enhance the growth and aerotolerance of Campylobacter spp. The main disadvantage of CCDA medium is its dark opacity, making it difficult to differentiate between Campylobacter spp. and non-Campylobacter spp. flora. Cefoperazone does not inhibit growth of molds and yeast on CCDA medium which can be associated with poultry samples.\nRothenberg et al (Applied and Environmental Microbiology, Volume 48(1), 78-80, Jul. 1984) disclose an attempt to develop an enrichment broth requiring only 7 hours of incubation and the comparison of their broth with that described by Doyle and Roman (Applied Environmental Microbiology, Volume 43, 1343-1353, 1982) and Park and Stankiewicz (Abstr. Assoc. Of Anal. Chem., Annu. Meet., volume 19, page 3, 1982). The Rothenberg et al. medium was a modification of the Doyle Roman broth and additionally contains 0.2% ferrous sulfate, 0.025% sodium metabisulfite, 0.05% sodium pyruvate, 0.1% sodium lauryl sulfate, and 0.075% agar. The Doyle and Roman broth contains Brucella broth, 7% lysed horse blood, 0.3% sodium succinate, 0.01% cysteine hydrochloride, vancomycin (15 xcexcg/ml), trimethoprim (5 xcexcg/ml), polymyxin B (20 IU/ml), and cycloheximide (50 xcexcg/ml). After a 16 to 18 hour incubation, the medium containing inoculum is plated directly onto Campy-BAP agar plates. The Stankiewicz broth contains vancomycin (20 mg/l), trimethoprim (10 mg/l), polymyxin B (5,000 IU/l for monophasic broth and 7,500 IU/l for diphasic medium), and lysed horse blood (5%, optional) in brucella broth. For the monophasic medium, 50 ml of the enrichment broth is placed in a 250 ml-Erlenmeyer flask. For the diphasic medium, brucella agar base (30 ml) with an overlay of the enrichment broth (50 ml) was made in a 500 ml-Erlenmeyer flask. Two other media are disclosed by Rothenberg. One contains 25 ml Brucella broth supplemented with 0.2% ferrous sulfate, 0.025% sodium metabisulfite, and 0.05% sodium pyruvate. The other, a selective medium, contains blood agar base no. 2 (Oxoid Ltd., London, England), 5% lysed horse blood, 10 xcexcg/ml vancomycin, 2.5 IU/ml polymyxin, 5 xcexcg/ml trimethoprim, and 15 xcexcg/ml cephalothin.\nCastillo-Ayala (Journal of Food Protection, Volume 55(5), 333-336, May 1992) disclose enrichment broths for isolation of Campylobacter jejuni/coli from freshly deboned market chicken. The first broth, VTP broth, contained 25 ml of double strength Brucella broth, 0.025% of sodium metabisulfate, 0.025% sodium pyruvate, 20 xcexcg/ml vancomycin, 10 xcexcg/ml trimethoprim, 5 IU polymyxin B. The other, BCN, contained 25 ml of double strength Brucella broth, 0.025% sodium metabisulfate, 0.025% sodium pyruvate, 50 IU/ml bacitracin, 20 IU/ml novobiocin, and 10 xcexcg/ml cycloheximide. They concluded that the vancomycin-trimethoprim-polymyxin B mixture was not a suitable agent for use in an enrichment-plating procedure to recover Campylobacter from poultry.\nChristopher et al (Journal of Food Protection, Volume 45(3), 260-262, Feb., 1982) disclose a method and media for isolation and enumeration of Campylobacter fetus subsp. jejuni that included a subculturing step that used Brucella broth containing 0.15% agar, 0.05% sodium pyruvate, 10 mg/liter vancomycin, 5 mg/liter trimethoprim, 2,500 IU/liter polymyxin B sulfate, 2 mg/l amphotericin B, and 15 mg/l cephalothin followed by subsequent streaking on plates of Brucella agar supplemented with 10% defibrinated horse blood and 10 mg/liter vancomycin, 5 mg/liter trimethoprim, 2,500 IU/liter polymyxin B sulfate, 2 mg/liter amphotericin B, and 15 mg/liter cephalothin.\nU.S. Pat. No. 5,302,388 (Doyle et al; Apr. 12, 1994) discloses an enrichment broth for Campylobacter jejuni containing Brucella broth, 7% lysed horse blood, 0.3% sodium succinate or 0.01% cysteine hydrochloride, 15 xcexcg vancomycin/ml, 5 xcexcg trimethoprim/ml, 20 IU polymyxin B/ml, and 50 xcexcg cycloheximide/ml.\nLuechtefeld and Wang (J. Clin. Microbiol., Volume 15(1), 137-140, January 1982) disclose the use of 1 mg or 400 xcexcg of 2,3,5-triphenyltetrazolium chloride (TTC) per ml of a medium containing brucella agar and sheep blood for differentiating C. fetus subsp. jejuni from C. fetus subsp. intestinalis. They report that all strains of C. fetus subsp. intestinalis were sensitive to 400 xcexcg/ml of TTC and all strains of C. fetus subsp. jejuni were insensitive to 400 xcexcg/ml TTC. The reference also discloses that Veron and Chatelain (Intl. J. Syst. Bacteriol., Volume 23, 122-134, 1973) found that none of the 18 strains of C. fetus subsp. fetus (C. fetus subsp intestinalis of Smibert (supra)) grew on blood agar containing 1 mg/ml of TTC in contrast to 9 strains of Campylobacter coli and 1 strain of C. jejuni. \nHanninen (Acta vet. scand., Volume 23, 88-98, 1982) disclose testing the tolerance of the C. jejuni/coli group to 1 mg/liter of TTC. The reference states that most strains in their study tolerated TTC. It further states that according to Skirrow and Benjamin, C. jejuni strains are sensitive to TTC and C. coli strains tolerated it. The reference further states that there are differences in TTC tolerance between the Campylobacter strains investigated by different authors and that most of the TTC tests made earlier and in their studies have been performed on a blood agar substrate containing TTC. Hanninen concludes that the blood in the substrate may be one reason why so many of the present strains tolerated TTC and the reference also observed that Skirrow and Benjamin did not use a blood-containing medium in their tests. The reference further discloses that the optimum TTC concentration for TTC reduction of C. fetus and related vibrios is shown to be 400 xcexcg/ml and notes that TTC has been used in a concentration range of 400 xcexcg/ml to 1 mg/ml.\nExisting media for the recovery and enumeration of colonies of Campylobacter species often contain many colonies of contaminating microorganisms and are difficult to enumerate because of the translucent nature of Campylobacter colonies. Therefore, there is a need in the art for media which recover high populations of Campylobacter with fewer contaminating colonies and which are easily enumerated. The present invention, described below, is improved media which provide selectivity for Campylobacter species while increasing the contrast of the colonies with the media to simplify counting procedures and is different from related art media.\nIt is therefore an object of the present invention to provide improved media for recovery and enumeration of Campylobacter species.\nAnother object of the present invention is to provide improved media for the recovery and enumeration of Campylobacter species, the improvement includes an indicator for facilitating accurate enumeration of Campylobacter colonies.\nA further object of the present invention is to provide improved media for the recovery and enumeration of Campylobacter species wherein the indicator is 2,3,5 Triphenyltetrazolium chloride (TTC).\nA still further object of the present invention is to provide improved media for the recovery and enumeration of Campylobacter species wherein the 2,3,5 Triphenyltetrazolium chloride concentration is about 200 xcexcg/ml or less.\nA further object of the present invention is to provide improved media for the recovery and enumeration of Campylobacter species which is a blood-containing or a blood-free media.\nA still further object of the present invention is to provide a media which contains a selective agent composition to provide improved selectivity for Campylobacter species.\nFurther objects and advantages of the present invention will become apparent from the following description.\nThe media of this invention may be used for the isolation and enumeration of Campylobacter species such as, for example, C. jejuni, C. coli, C. lari, etc., from a variety of sources. Although the media are particularly valuable for the growth and recovery of this microorganism from samples taken from poultry carcasses, especially chicken, it is understood that the media may also be used for the isolation of Campylobacter from any samples suspected of containing this pathogen. Without being limited thereto, other sources include animal carcasses such as cattle and sheep, food, milk, water, or environmental samples, or clinical sources such as blood or feces. The culture media contemplated for use in this invention may be prepared using techniques conventional in the art. The basal medium components including agar and/or nutrient media with an energy source are mixed, heated to boiling and sterilized by autoclaving. After cooling the sterilized medium to about 50xc2x0 C.-55xc2x0 C., blood and filter-sterilized supplements are added with mixing, pH asceptically adjusted to about pH 7.4 with, for example, 10N NaOH, and the medium finally poured into a culture container, such as a petri dish, for example, and cooled to allow the agar to solidify.\nIn one embodiment of the present invention, the basal medium components selected for use are not critical and may be readily determined by the practitioner skilled in the art, one of the improvements includes the addition of an indicator, 2,3,5-triphenyltetrazolium chloride in amounts which do not significantly inhibit the growth of Campylobacter species. Any nutrient medium and energy source effective to support growth of Campylobacter species may be used. Suitable nutrient media include, for example, but are not limited to Brucella agar (e.g. BBL, Cockeysville, Md.; Difco Laboratories, Detroit, Mich.;) CM 691 (Oxoid, Columbia, Md.), Campylobacter agar base (Difco), Blood agar base No. 2 (Oxoid), Brain-heart infusion agar (BBL;Difco), or Columbia Blood Agar Base. A variety of energy sources may also be employed, and may be incorporated into commercially available nutrient media or added separately. In the preferred embodiments of the present invention, Brucella agar (Acumedia) is the preferred nutrient medium. Suitable energy sources for use in the media are described by Sibert (INL Bergey\"\"s Manual; supra, the contents of which are herein incorporated by reference) and include pyruvate, citrate, succinate, cis-aconitate, xcex1-ketoglutarate, fumarate, malate, and oxaloacetate.\nFor all embodiments of the present invention, the source of blood added to the medium also is not critical. While horse blood is preferred, it is understood that other blood sources may be used, such as for example, sheep blood. For blood-free media, hemin, yeast extract, sodium carbonate and xcex1-ketoglutaric acid are used to replace the blood component.\nSelective agents, for use in all embodiments of the present invention, are used to prevent the growth of contaminating microorganisms present in samples to be tested, but which do not inhibit growth of Campylobacter species. The selective agents include any selective agent or combination thereof known to one of ordinary skill in the art. Improved preferred selective agents include vancomycin, trimethoprim, polymyxin B, cycloheximide, rifampicin, nystatin, amphotericin, and cefoperazone, their salts, and mixtures thereof. The combination of these preferred agents provides excellent selectivity for Campylobacter species.\nAn indicator, for use in all embodiments of the present invention, is included in Campylobacter media to facilitate manual or automated enumeration of Campylobacter colonies. The indicator 2,3,5-triphenyltetrazolium chloride (TTC; Sigma) is used in a concentration range of about 50 xcexcg/ml to about 250 xcexcg/ml. The preferred concentration is about 200 xcexcg/ml.\nOther adjuvants, useful in all embodiments of the present invention, may also be incorporated into the media for enhancing growth and/or aerotolerance of Campylobacter. Preferred adjuvants enhancing aerotolerance are described by Smibert (supra) and include but are not limited to sodium pyruvate or pyruvic acid, ferrous sulfate, bovine superoxide dismutase, catalase and reducing agents such as sodium bisulfite or sodium metabisulfite. Particularly preferred for addition to medium are ferrous sulfate, pyruvic acid, and sodium bisulfite. It is understood that the use of blood in the media also enhances aerotolerance because it contains heme, catalase and superoxide dismutase.\nThe concentration and amount of each of the components of the Campy-TTC blood-containing medium, a particularly preferred embodiment of the present invention, are variable and may be readily determined by the practitioner skilled in the art. The amount of each component of the basal or nutrient media should be effective to promote growth of Campylobacter species, while the amount of the selective agents should be effective to inhibit growth of contaminating (non-Campylobacter) microorganisms without substantially inhibiting growth of Campylobacter species relative to culture medium lacking these selective agents. Without being limited thereto, preferred ranges of the selective agents include about 20-50 mg/liter cefoperazone, about 100-400 mg/liter cycloheximide, about 5-20 milligrams/liter vancomycin, about 2.5-10 milligrams/liter trimethoprim, about 5-20 milligrams/liter polymyxin, about 5-20 milligrams/liter rifampicin, about 10-40 milligrams/liter amphotericin, about 20-80 milligrams/liter nystatin, and about 100-250 mg/liter 2,3,5-Triphenyltetrazolium chloride.\nIn accordance with a particularly preferred formulation of blood-containing Campy-TTC, ranges of the amount of each component per liter include but are not limited to:\nThe final pH of the medium should generally be between about 6.5 to about 7.8, pH 7.4 is preferred.\nIn accordance with another particularly preferred formulation, blood-free Campy-TTC, ranges of the amount of each component per liter include but are not limited to:\nThe final pH of the medium should generally be between about 6.5 to about 7.8, pH 7.4 is preferred.\nIn accordance with another particularly preferred formulation, blood-free Campy-TTC, ranges of the amount of each component per liter include but are not limited to:\nfinal pH of the medium should generally be between about 6.5 to 7.8, pH 7.4 is preferred.\nIn accordance with another particularly preferred formulation, blood-free Campy-TTC, ranges of the amount of each component per liter include but are not limited to:\nThe final pH of the medium should generally be between about 6.5 to 7.8, pH 7.4 is preferred. Effective amounts of each ingredient are those amounts which promote growth of Campylobacter species.\nIn accordance with another particularly preferred formulation, blood-free Campy-TTC, ranges of the amount of each component per liter include but are not limited to:\nThe pH of the medium should generally be between about 6.5 to 7.8, pH 7.4 is preferred. Effective amounts of each ingredient are those amounts which promote growth of Campylobacter species.\nIn accordance with another particularly preferred formulation, blood-free campy-TTC, ranges of the amount of each component per liter include but are not limited to:\nThe final pH of the medium should generally be between about 6.5 to 7.8, pH 7.4 is preferred. Effective amounts of each ingredient are those amounts which promote growth of campylobacter species.\nIn use, the sample to be analyzed is inoculated onto the culture medium using techniques conventional in the art and is incubated for a sufficient time and under conditions effective to promote growth of Campylobacter species. Suitable conditions may be readily determined by the practitioner skilled in the art and are described by Smibert (supra). Without being limited thereto, preferred conditions include a temperature between about 35xc2x0 C. to about 44xc2x0 C., especially about 42xc2x0 C. to about 43xc2x0 C., and a low oxygen tension (i.e., microaerobic), especially an oxygen concentration of between about 3-6%. Techniques for generating this reduced atmosphere are well known in the art and are described in, for example, Inoue (U.S. Pat. No. 4,904,597), or Hutchinson and Bolton (J. Clin. Pathol., Volume 36, 1350-1352, 1983), or Park et al. (supra), the contents of each of which are herein incorporated by reference.\nFollowing incubation, generally after about 24-48 hours, the culture may be examined for the presence of colonies indicative of Campylobacter. Colonies of Campylobacter on any of the preferred Campy-TCC media of this invention are deep red to magenta in color. These colonies can be readily discriminated from non-Campylobacter breakthrough flora. Further confirmatory testing of the colonies and speciation can be conducted as described by Park et al (supra) or Smibert (supra). The use of these agars containing 2,3,5-triphenyltetrazolium chloride would be a benefit to technicians in enumerating and confirming Campylobacter colonies. The use of Campy-TTC agars, either blood-containing or blood-free with a preferred selective agent composition, also is a benefit because of the increased selectivity of the agars. Automated counting methods may also be made possible using triphenyltetrazolium chloride containing agars since the darker colonies contrast sufficiently with the agar base to allow detection by electronic means."}
{"text": "The present invention relates to a gear assembly and, more particularly, to such a gear assembly which allows ease of adjustment of the relative position of its components. The invention is particularly applicable to gear assemblies for use in the steering of automotive vehicles.\nA gear assembly of the rack and pinion type comprises a toothed nut located on a threaded input shaft so that rotation of the input shaft produces an axial displacement of the nut. The teeth of the nut engage with corresponding pinion teeth on an output shaft so that displacement of the rack nut produces a rotation of the output shaft. In such gear assemblies it is desirable to be able to adjust the relative positions of the rack and the pinion to ensure optimum engagement of the teeth.\nThe document U.S. Pat. No. 4,614,127 describes a rack and pinion gear assembly in which the output shaft is journaled in the housing by way of two eccentric bushings. Rotation of a bushing in the housing produces an angular displacement of the output shaft. Thus, by selected rotation of the bushings, the output shaft can be displaced to bring its teeth into optimum engagement with those of the rack. Once the optimum position of the bushings has been reached they are rotationally locked by punching part of the bushing into a corresponding notch in the housing. While this device ensures an optimum interengagement of the teeth, the locking of the bushings by punching is not completely satisfactory as the punching operation can cause undesirable displacement of the bushing.\nThe object of the present invention is to provide a gear assembly which allows optimum alignment of the gear components and in which the locking of the bushings is both simple and efficient."}
{"text": "1. Field of the Invention\nThe present invention relates generally to transaction systems and methods for verifying a consumer engaged in a transaction with a merchant in order to authorize both the consumer as well as the transaction and, in particular, to a computer-implemented method, system and apparatus for dynamically verifying a consumer engaged in a transaction with a merchant and authorization of this transaction.\n2. Description of Related Art\nIn order to enable convenient purchases of goods and services by consumers, the financial service industry has developed many alternative payment methods that allow a consumer to engage in a transaction and receive goods and services on credit. For example, such alternative payment methods may include checks, ATM or debit cards, credit cards, charge cards, etc. prior to the birth of virtual commerce, as discussed below, such payment options provide an adequate convenience and transactional security to consumers and merchants in the market place. While transactional security may include the security offered by a payment method to the consumer that the purchase event will not result in a breach of personal information or that the consumer is a victim of identity theft, transactional security also offers the merchant or seller the security that fraud will not be perpetrated.\nVirtual commerce and the growth of the Internet as a medium for commerce have placed pressure on the payment options discussed above on both the convenience and transactional security by the credit issuer. For example, credit cards may be convenient to the consumer, but are subject to fraudulent use via theft of the account number, expiration date and address of the consumer. This, in turn, places the credit issuer at risk of offering credit to an uncreditworthy consumer, being the subject of consumer fraud or providing authorization to a merchant to provide services or ship goods to a fraudulent source.\nCurrently available payment options include significant shortcomings when applied to remote purchasers, such as purchases where the buyer and the seller (that is, the merchant) are not physically proximate during the transaction. Further, regardless of the proximity of the consumer and the merchant, merchants and credit issuers alike continue to battle the problem of fraudulent purchases. Each new payment option and every new sales channel (in-store, telephone, mail and Internet) have, in turn, spawned innovation on the part of consumers willing to perpetrate fraud in order to obtain goods and services without paying for them.\nIn recent years, the birth of the Internet commerce industry and the continued growth in mail order and telephone order commerce have pushed the credit card to the forefront of these battles. Typically, merchants are forced to rely on credit cards because it is currently their only option in the remote purchase environment. However, regardless of the type of credit offered, low transactional security is offered to both merchants and consumers. This leads to significant cost for the consumers and the merchants, such as the consumer cost including the impairment of their credit record, the inconvenience of changing all of their credit card accounts and the financial cost in resolving the situation. Merchant costs may include the mitigation of fraud losses, including the cost in incremental labor, hardware and software to implement additional security checks in their sales/order entry software, higher transaction processing expense in the form of discount rates for credit cards and NSF fees for checks and higher fraud charge-offs for undetected fraudulent purchases.\nWith the continuing speed and ability of a consumer to gain credit, whether at a point-of-sale or through the use of an existing account, identity theft and fraud are on the increase. Fortunately, those that perpetrate this fraud and theft typically target specific consumers, such as the elderly, and engage in specified and easily-recognizable patterns, such as elevated purchase costs at otherwise lower-end shopping facilities in a compressed period of time. Further, many fraud perpetrators and thieves have a specific pattern of buying that is easily assessed by a third party.\nPresently, merchants attempt to properly identify a consumer using a credit card or other credit account at the point-of-sale. However, as is easily evident in today's marketplace, merchants are often more interested in providing a consumer with quick and efficient service with little hassle regarding the consumer's identity. Further, merchants do not have the data available to them at the point-of-sale for making an appropriate identification of a consumer or otherwise detecting a fraudulent purchase. While some merchants do use external databases to verify a consumer and authorize the transaction, these databases include errors, have limited information, have data omissions and further include data that may be compromised by an external source. Therefore, there remains a need for a more dynamic verification process for both verifying the consumer and authorizing the transaction prior to final acquisition of the goods and/or services by the consumer."}
{"text": "Redeye is the appearance of an unnatural reddish coloration of the pupils of a person appearing in an image captured by a camera with flash illumination. Redeye is caused by light from the flash reflecting off blood vessels in the person's retina and returning to the camera.\nA large number of image processing techniques have been proposed to detect and correct redeye in color images. In general, these techniques typically are semi-automatic or automatic. Semi-automatic redeye detection techniques rely on human input. For example, in some semi-automatic redeye reduction systems, a user must manually identify to the system the areas of an image containing redeye before the defects can be corrected.\nMany automatic redeye reduction systems rely on a preliminary face detection step before redeye areas are detected. A common automatic approach involves detecting faces in an image and, subsequently, detecting eyes within each detected face. After the eyes are located, redeye is identified based on shape, coloration, and brightness of image areas corresponding to the detected eye locations. In general, face-detection-based automatic redeye reduction techniques have high computation and memory resource requirements. In addition, most of the face detection algorithms are only able to detect faces that are oriented in an upright frontal view; these approaches cannot detect faces that are rotated in-plane or out-of-plane with respect to the image plane.\nA typical prior art redeye filter process is illustrated in FIG. 1(a). An input image is first analyzed by a speed optimized redeye detection stage 100 at a pixel level 103 and segmented into candidate redeye regions 104. A further series of falsing and verification filters 106 are then applied to the candidate regions and a set of confirmed redeye regions 108 is thus determined. A correction filter (pixel modifier) 102 is next applied to the confirmed regions and a final image 112, corrected for redeye, is generated.\nExemplary prior art includes U.S. Pat. No. 6,407,777 to DeLuca which discloses in-camera detection and correction of redeye pixels in an acquired digital image; U.S. patent application 2002/0176623 to Steinberg which discloses automated real-time detection and correction of redeye defects optimized for handheld devices; U.S. patent applications 2005/0047655 and 2005/0047656 to Luo et al which disclose detecting and correcting redeye in a digital image and in embedded systems respectively.\nNow it is well known that within an image acquisition subsystem such as is embodied in typical digital cameras, the peak computing load and resource requirements occur around the time of image acquisition. Upon receiving an image acquisition request from the user the main embedded processing system must refine the image focus and exposure to achieve an optimal main acquired image; this image, in turn, must be off-loaded from the main optical sensor of the camera and subjected to further image processing to convert it from its raw format (e.g. Bayer) to a conventional color space such as RGB or YCC. Finally the acquired image must be compressed prior to saving it on a removable storage medium such as a compact flash or multimedia card.\nThe time taken by the camera to recover from the acquisition of a first image and reinitialize itself to capture a second image is known in the industry as the “click-to-click” time. As this is one of the most important parameters for the comparison and marketing of modern digital cameras it vital for manufacturers to minimize said “click-to-click” time. Thus any additional image processing, such as redeye filtering, which is to be added to the main image acquisition chain should be highly optimized for speed of execution in order to minimize its impact on the click-to-click time of the main system.\nEvidently such a redeye filter must compromise its overall performance in terms of accuracy of detection of redeye defects and quality of image correction. An alternative would be to wait until after the main image has been acquired and perform the redeye filtering at a later time when the camera may execute the filter as a background process, or to perform the redeye filtering off-camera on a desktop PC or printer.\nHowever there are some drawbacks to this approach. Firstly, images will be displayed on the acquiring device, immediately after acquisition, with uncorrected redeye defects; and, when images are accessed in playback mode, there will be a further delay while images are post-processed before an image can be displayed. Both drawbacks would create a negative impression on end users.\nFurther, as practically all digital cameras store images using lossy compression techniques there are additional disadvantages with respect to image quality as images must be decompressed and recompressed in order to perform the redeye detection and correction processes in playback or background modes. Such loss of image quality may not become apparent until later when a user wishes to print an image and it is too late to reverse the process.\nIf redeye processing is delayed until the images are loaded onto another device, such as a desktop PC or printer there are further disadvantages. Firstly, important meta-data relating to the acquiring device and its state at the time the image was acquired may not be available to the redeye filter process. A second disadvantage is that this post-processing device must perform redeye filtering on the entire image; where this is an embedded device such as a printer it may, itself, be relatively constrained in terms of CPU cycles and processing resources for its primary post-processing activity and it may be desirable to optimize the performance of the full redeye filter."}
{"text": "In the field of “Oil & Gas”, reciprocating compressors are widely used.\nUS2002/0141884-A1 describes an unloader system is provided for a reciprocating gas compressor. The system includes an unloader valve assembly including a valve member controlling flow between compressor cylinder and a clearance bottle. Opening and closing of the valve member is controlled by manipulating a control pressure acting through a manifold against the stem of the valve member by means of a pressure regulator connected in series with a pressure source. When the pressure in the compressor cylinder acting on the heads of the poppet valve members exceeds the control pressure acting on the stems, the poppet valve members open, partially unloading the compressor.\nEach cylinder head of the reciprocating compressor described in the cited document has an annular chamfer with a concave surface partially facing the valve openings.\nIn different known designs, the cylinder head may a have plurality of slots, each facing a valve opening. Each slot have a concave surface, shaped as a part of a sphere (therefore with a constant radius).\nThe fluid, when aspired into the compression chamber and when leaving it, runs onto the concave surface of the chamfer or of the slots. This generates pressure losses that lower the efficiency of the compressor."}
{"text": "Implantable electrical stimulation systems have proven therapeutic in a variety of diseases and disorders. For example, spinal cord stimulation systems have been used as a therapeutic modality for the treatment of chronic pain syndromes. Deep brain stimulation has also been useful for treating refractory chronic pain syndromes and has been applied to treat movement disorders and epilepsy. Peripheral nerve stimulation has been used to treat chronic pain syndrome and incontinence, with a number of other applications under investigation. Functional electrical stimulation systems have been applied to restore some functionality to paralyzed extremities in spinal cord injury patients. Moreover, electrical stimulation systems can be implanted subcutaneously to stimulate subcutaneous tissue including subcutaneous nerves such as the occipital nerve.\nStimulators have been developed to provide therapy for a variety of treatments. A stimulator can include a control module (with a pulse generator), one or more leads, and an array of stimulator electrodes on each lead. The stimulator electrodes are in contact with or near the nerves, muscles, or other tissue to be stimulated. The pulse generator in the control module generates electrical pulses that are delivered by the electrodes to body tissue."}
{"text": "1. Field of the Invention\nThe present general inventive concept relates to a carrier recovery apparatus usable with a vestigial side band (VSB) receiver and a method thereof. More specifically, the present general inventive concept relates to a carrier recovery apparatus using not only a pilot signal, but also upper and lower sidebands of a received VSB signal, and a method thereof.\n2. Description of the Related Art\nIn order for a data receiver to accurately demodulate VSB-modulated data, it is necessary to minimize frequency offset and high levels of phase noise (jitter) generated from a tuner or an RF (radio frequency) oscillator used for data receiving. This procedure is called ‘carrier recovery’.\nA digital broadcast system based on a VSB modulation method in conformity with the standards of the Advanced Television Systems Committee (ATSC) uses a pilot signal in a transmitting signal for carrier synchronization. Here, the pilot signal is a signal loaded on a carrier during the transmission for accurate carrier recovery.\nFIG. 1 is a block diagram illustrating a conventional carrier recovery circuit 100 for phase detection. Referring to FIG. 1, the conventional carrier recovery circuit 100 comprises a multiplier 101, a pilot detector 103, a phase detector 105, a loop filter 107, and a numerically controlled oscillator (NCO) 109.\nA received VSB signal is digitalized by an analog-to-digital converter (ADC) (not shown), and output as a baseband signal by the multiplier 101.\nThe pilot detector 103 detects a pilot signal from the baseband signal, and the phase detector 105 reads phase information of the pilot signal. There are many methods for reading phase information of the pilot signal, and a suitable method is selected depending on the application. The phase information read by the phase detector 105 goes through the loop filter 107, is converted to a frequency component through the NCO 109, and is multiplied with the received VSB signal by the multiplier 101 to output the baseband signal.\nThe above-described procedure is repeatedly performed according to a feedback operation until a phase error of the pilot signal becomes zero.\nHowever, if the pilot signal in the wireless environment is corrupted, a received signal cannot be recovered properly. The unwanted pilot signal corruption occurs often in a multi-path environment, and eventually causes performance degradation of the VSB receiver. This problem can occur in any American digital broadcasting system using the VSB method."}
{"text": "Latex binder compositions have been used in a wide variety of applications such as adhering paper coatings to cellulosic substrates, in carpet and rug backings, and in various non-woven products which have gained broad acceptance as replacements for woven fabrics, including such articles as facings or top sheets in diapers, incontinence pads, sanitary napkins, hospital gowns, and other single and multiple-use non-woven products.\nConventional latex binders have generally been prepared from acrylic polymers, styrene/acrylate copolymers, or styrene/butadiene copolymers containing N-methylol functionality which, upon curing, generate substantial quantities of formaldehyde, typically from 200 to 500 ppm or more. It is desirable to use binders which are essentially free of formaldehyde generation. Formaldehyde-free latex binders include various urethane polymers, acrylic polymers, styrene/acrylate copolymers, and the like, such as those disclosed in U.S. Pat. Nos. 2,837,462; 4,207,367; 4,381,332; 5,030,507, and others. However, none of the prior art discloses a suitable conjugated diene/vinyl-substituted aromatic copolymer latex binder which is free of formaldehyde formation, and which also exhibits superior coagulation stability and mechanical stability while imparting excellent wet tensile strength to cellulosic articles such as paper."}
{"text": "In the effort toward achieving higher data density on recording media, spin filters have become of interest for use in magnetoresistive (MR) read heads. FIG. 1 is a diagram of a conventional spin filter 10. In general, the conventional spin filter 10 would be incorporated into a MR read head (not explicitly shown), which would include leads electrically connected to other electronics to drive current through the conventional spin filter 10 during reading. In such an application, current is generally driven in the current perpendicular to the plane (CPP) configuration. The CPP configuration is in the z-direction depicted in FIG. 1.\nThe conventional spin filter 10 includes a seed layer 12, an antiferromagnetic (AFM) layer 14, a pinned layer 16, a nonmagnetic spacer layer 24, a free layer 26, a filter layer 28, a specular oxide layer 30, and a capping layer 32. The seed layer 12 is used to provide the appropriate surface for growing the AFM layer 14 with the desired crystal structure. The AFM layer 14 is used in pinning the magnetization of the pinned layer 16. The pinned layer 16 may be a synthetic pinned layer, including ferromagnetic layers 18 and 22 separated by an electrically conductive spacer layer 20 that is typically Ru. The electrically conductive spacer layer 20 has a thickness configured to ensure that the ferromagnetic layers 18 and 22 are antiferromagnetically coupled. Thus, the magnetization of the ferromagnetic layer 18 is pinned by the AFM layer 14. The magnetization of the ferromagnetic layer 22 is set because it is strongly antiferromagnetically coupled to the magnetization of the ferromagnetic layer 18. The nonmagnetic spacer layer 24 is typically electrically conductive, for example Cu. The free layer 26 is ferromagnetic and typically includes materials such as CoFe.\nThe filter layer 28 has a high electrical conductivity and typically includes materials such as Cu. The specular oxide layer 30 may be a nano-oxide and typically includes materials such as alumina. The combination of the filter layer 28 and the specular oxide layer 30 has been used to provide adequate specularity of scattering of electrons from the free layer 26 that are incident on the specular oxide layer 30. In particular, the filter layer 28 has been utilized in the conventional spin filter 10 to provide a region, or filter, for specular reflection between the free layer 26 and specular layer 30. Thus, the magnetoresistance of the conventional spin filter 10 is improved. Consequently, the magnetoresistance of the conventional spin filter 10 is adequate. The capping layer 32 is typically oxidized Ta.\nAlthough the conventional spin filter 10 functions, there are drawbacks to the use of the conventional spin filter 10. Insertion of the specular oxide layer 30 can increase the coercivity of the free layer 26, which is undesirable. Furthermore, the specular oxide layer 30 is generally a nano-oxide that can continue to oxidize during processing. The signal may degrade during the lifetime of the conventional spin filter 10. The conventional spin filter 10 thus suffers thermal instabilities and may have reduced reliability.\nAnalogous conventional spin filters are described in U.S. Pat. No. 6,795,279 B2; U.S. Pat. No. 6,556,390 B1; U.S. Pat. No. 5,898,612; U.S. Pat. No. 6,407,890 B1; U.S. Pat. No. 6,764,778 B2; U.S. Pat. No. 6,700,753 B2; U.S. Pat. No. 6,775,111 B2; U.S. Pat. No. 6,591,481; U.S. Pat. No. 6,613,380 B1; U.S. Pat. No. 6,636,398 B2.\nAccordingly, what is needed is a system and method for providing a spin filter having improved thermal stability, signal sensitivity, and/or reliability."}
{"text": "U.s. pat. No. 2,953,460 PA1 U.s. pat. No. 3,708,255 PA1 U.s. pat. No. 3,882,768"}
{"text": "An intraluminary prosthesis is a medical device used in the treatment of diseased bodily lumens. One type of intraluminary prosthesis used in the repair and/or treatment of diseases in various body vessels is a stent. A stent is generally a longitudinal tubular device formed of biocompatible material which is useful to open and support various lumens in the body. For example, stents may be used in the bodily vessel, such as in the coronary or peripheral vasculature, esophagus, trachea, bronchi colon, biliary tract, urinary tract, prostate, brain, as well as in a variety of other applications in the body. These devices are implanted within the vessel to open and/or reinforce collapsing or partially occluded sections of the lumen.\nStents generally include an open lattice configuration. This configuration allows the stent to be inserted through curved vessels. Furthermore, this configuration allows the stent to be configured in a radially compressed state for intraluminary catheter implantation. Once properly positioned adjacent the damaged vessel, the stent is radially expanded so as to support and reinforce the vessel. Radial expansion of the stent may be accomplished by inflation of a balloon attached to the catheter or the stent may be of the self-expanding variety which will radially expand once deployed. Tubular shaped structures, which have been used as intraluminary vascular stents, have included helically wound coils which may have undulations or zig-zags therein, slotted stents, ring stents, braided stents and open mesh wire stents, to name a few. Super-elastic materials and metallic shape memory materials have also been used to form stents.\nWhile stents are often made from metallic materials, the use of plastic stents is not uncommon, especially in non-vascular applications. For example, plastic stents have been used to treat malignant or benign strictures throughout the gastrointestinal tract because of, among other things, ease of placement and non-permanency of the stents. In the case of biliary applications, re-intervention of a plastic stent after three months is typical for malignant tumors where after the plastic stent is replaced with another plastic stent or perhaps a metallic stent. Benign strictures in biliary applications are often treated every three months with a plastic stent for up to about a year. In duodenal applications, a plastic stent is often placed as a bridge to surgery. In esophageal applications, a plastic stent may be placed as an adjunct treatment to radiation.\nWhile plastic stents offer many advantages, a drawback of a typical plastic stent is its patency rate as compared to a metallic stent. A metallic stent generally has much larger radial force and a longer patency rate, especially when the stent is covered. Metal stents, however, have generally not been used in benign applications due to the difficulty of removing an implanted metallic stent.\nIn biliary applications, re-intervention may be required in about twenty-five percent of the time due to tumor in-growth through the open lattice portions of the implanted stent. Covering the stent, for example with a layer of silicone, may reduce the re-intervention due to tumor in-growth less than two percent of the time.\nU.S. Pat. No. 5,603,722 to Phan et al describes a stent made from plastic, more particularly a shape memory polymer. The stent includes elongate strips of shape memory polymer overlapping wound to form a closed-wall stent structure in a contacted state of the stent. When the stent is expanded portions of the strip segments separate from each other to provide a stent with an open lattice structure. U.S. Pat. No. 5,163,952 to Froix describes a stent made from plastic, more particularly a shape memory polymer, which is in the shape of a coil or a solid-walled tubular cylinder. While the coil is a flexible structure, it has gaps into which tumor in-growth may occur. The solid-walled tubular cylinder may avoid the tumor in-growth concerns associated with open lattice stent structures, but such a stent is not highly flexible, making delivery and placement into curved lumens difficult.\nThus, there is a need for a plastic stent which has improved patency by reducing re-intervention rates due, for example, to tumor in-growth, while still being flexible so that it can used in curved lumens."}
{"text": "1. Field of the Invention\nThe invention relates to an information carrier and processes for producing an original copy of the carrier, which comprises a recording layer deposited on a carrier material. A relief image containing the information is embossed in the recording layer.\n2. Description of the Prior Art\nImages with grating-like screening have been produced by the ZOD (zero order diffraction) technique described in the LASER u. Elektro-Optik Journal No. 3/1976, pages 16-17. Three nickel matrices are produced from the relief images which correspond, for example, to three primary-color grating patterns in a photoresist, and colorless thermoplastic films of, for example, polyvinyl chloride are embossed with these matrices. These films are mechanically superimposed, and, on projection with conventional projectors, colored projection images are obtained from the colorless relief images. The grating-shaped screening is effected with relief gratings of rectangular cross-section, the grating periods being approximately 1.5 .mu.m. A separate nickel matrix with different relief depths is made for each color separation (red, yellow and blue), from which separate embossed images are produced. The relief depths differ depending on the color separation. The greatest relief depths are used for the red color separation while the smallest are used for the blue color separation. The color separation images are screened. The embossed images are then superimposed to form a three-layered relief image which can be used to project colored images. The technique described yields very bright color images with high resolution. The relief images can be duplicated relatively cheaply and rapidly by an embossing process.\nA disadvantage, which has hindered the introduction of this technique, is the expansive production process, incurred by performing three completely separate operations for producing the individual, embossed relief images, corresponding to the color separations. A further disadvantage ensues from the necessity of aligning the three separate relief images to form the duplicate image required for the colored projection."}
{"text": "The present invention relates to fluids and compositions useful in subterranean operations, and more particularly, to additives and treatment fluids with improved shale inhibition, and associated methods of use.\nA treatment fluid may be used in a variety of subterranean operations. As used herein, the term “subterranean operation” is defined to mean any operation that requires the performance of some action or procedure below the surface of the earth, including, but not limited to, actions or procedures performed in the course of recovering oil, gas, and/or other substances from a formation below the surface of the earth. Such subterranean operations also include, but are not limited to, service line and tunneling operations. As used herein, the term “treatment,” or “treating,” does not imply any particular action by the fluid or any particular component thereof, but instead refers to any use related to a subterranean operation in conjunction with a desired function and/or for a desired purpose. For example, a fluid may be used to drill a well bore in a subterranean formation or to complete a well bore in a subterranean formation, as well as numerous other purposes.\nDrilling operations may involve any number of different techniques known and practiced in the art. In the most basic sense, rotary drilling operations typically involve attaching a drill bit on a lower end of a drillstring to form a drilling tool and rotating the drill bit along with the drillstring into a subterranean formation to create a well bore through which subsurface formation fluids may be recovered. In another method of drilling, coiled tubing may be used instead of jointed pipe and the drill bit may be rotated using a downhole motor. The process of drilling a well bore typically requires the use of a drilling fluid. Drilling fluids are used, inter alia, to cool the drill bit, lubricate the rotating drill pipe to prevent it from sticking to the walls of the well bore, prevent blowouts by serving as a hydrostatic head to counteract the sudden entrance into the well bore of high pressure formation fluids, and remove drill cuttings from the well bore. A drilling fluid used in connection with drilling a well in a subterranean formation may be any fluid substance (gaseous or liquid) or a mixture of fluids and solids (such as solid suspensions, mixtures and emulsions of liquids, gases and solids).\nDuring drilling of subterranean well bores, it is not uncommon to encounter various strata that comprise reactive shales. As used herein, the term “shale” is defined to mean materials that may “swell,” or increase in volume, when exposed to water. Examples of these shales include certain types of clays (for example, bentonite). Reactive shales may be problematic during drilling operations because of, inter alia, their tendency to degrade when exposed to aqueous media such as aqueous-based drilling fluids. This degradation, of which swelling is one example, can result in undesirable drilling conditions and undesirable interference with the drilling fluid. For instance, the degradation of the shale may interfere with attempts to maintain the integrity of drilled cuttings traveling up the well bore until such time as the cuttings can be removed by solids control equipment located at the surface.\nShale disintegration also may impact “equivalent circulating density” (“ECD”). ECD may be affected by the solids content of the drilling fluid, which may increase if surface solids control equipment cannot remove shale from the drilling fluid. Plastic viscosity (an indicator of size and quantity of solids) is an important parameter that affects drilling rate. Maintenance of appropriate ECD is important in drilling a well bore where a narrow tolerance exists between the weight of the drilling fluid needed to control the formation pressure and the weight of the drilling fluid that will fracture the formation. In such circumstances, minimizing shale degradation may be desirable, inter alia, to control of the viscosity of the drilling fluid. Moreover, degradation of drilled cuttings prior to their removal at the surface may prolong drilling time because shale particles traveling up the well bore can break up into smaller and smaller particles, which can expose new surface area of the shale particles to the drilling fluid and lead to further absorption of water and degradation.\nShale degradation may substantially decrease the stability of the well bore, which may cause irregularities in the diameter of the well bore, e.g., the diameter of some portions of the well bore may be either smaller or greater than desired. In an extreme case, shale degradation may decrease the stability of the well bore to such an extent that the well bore collapses. Degradation of the shale also may interrupt circulation of the drilling fluid, cause greater friction between the drill string and the well bore, and/or cause the drill string to become stuck in the well bore. Accordingly, the complications associated with shale swelling during drilling may substantially increase the time and cost of drilling.\nA traditional method of inhibiting shale swelling during drilling to attempt to minimize such complications has been to use an oil-based drilling fluid as opposed to an aqueous-based drilling fluid. However, environmental regulations enacted by numerous countries have limited the use of oil-based drilling fluids.\nAnother conventional technique used to counteract the propensity of aqueous drilling fluids to interact with reactive shales in a formation involves adding a shale-inhibiting component to the aqueous drilling fluid. As used herein, the term “shale-inhibiting component” refers to a compound that demonstrates a propensity for inhibiting the tendency of shale to absorb water. Amphoteric materials (i.e., substances that may exhibit both acidic and/or alkaline properties) are one type of water-based shale inhibitor that has been used in the past. Amphoteric materials are believed to attach to the shale substrate, thus preventing water ingress. However, amphoteric inhibitors may be environmentally undesirable, especially in heavily regulated areas, because they typically demonstrate low biodegradability and high toxicity. Potassium chloride is another conventional shale-inhibiting component. Although potassium chloride is widely used as a shale inhibitor in the North Sea, it is considered to be only moderately effective at inhibiting shale swelling. Furthermore, potassium chloride can be environmentally unacceptable in other areas of the world, for example, regions such as around the Gulf of Mexico, because its concentration of potassium ions may harm certain types of marine life. Potassium chloride is also disfavored in regions such as the Middle East where wells are drilled in close proximity to aquifers due to concerns that the potassium chloride will contaminate the aquifer. Polyglycols have also been used as shale inhibitors in water-based drilling fluids, but have not demonstrated satisfactory inhibition levels. Partially hydrolyzed polyacrylamides (“PHPA”) have also been utilized in many regions, but these have been found to have undesirable properties in certain circumstances.\nDrilling fluids containing silicates are used in the art to manage shale reactivity with aqueous-based drilling fluids. However, silicates do not provide the desired lubricity (i.e., the ability to lubricate equipment used in a drilling operation) for subterranean drilling operations, nor do they effectively reduce the accumulation of shale on equipment that comes into contact with those shales encountered in a well bore. Drilling fluids containing nanoparticles are also used in the art to manage shale reactivity with aqueous-based drilling fluids by, inter alia, providing lubricity and preventing the accretion of shale on drilling equipment. However, these nanoparticles generally do not impart increased hardness to the shale, a property that may facilitate the removal of shale particles from a well bore and/or circulated drilling fluid."}
{"text": "1. Field\nThis invention relates to the hot mix asphalt plant used to manufacture asphaltic mixes to pave highways.\n2. Description of the Related Art\nThe components of the state of the art hot mix asphalt plant combine process methods to heat and dry aggregates, add coarse and fine aggregate dust particles and mix them with asphalt. The raw material virgin aggregates release dust particles as they are heated and dried that must be collected using mechanical separation for the coarse dust particles along with fine dust collection using fabric filters or baghouses to eliminate air borne particulate from entering the atmosphere. Asphalt plants use rotary dryers to tumble the aggregates in the presence of heat to remove the moisture and heat them to the process temperature.\nDrying aggregates with the rotary dryer involves a physical separation process designed for the purpose to remove the liquid phase or moisture using thermal energy. The liquid in this case is water that is liberated by the process of vaporization leaving a solid or aggregates with a trace of residual moisture. The dryer is a direct heat dryer as heat transfer for drying is accomplished by direct contact between the wet solid aggregate with hot gases from a combustion burner. This invention describes the dryer as a rotating cylindrical shell, slightly inclined to the horizontal with the heating medium flowing in a countercurrent direction to the flow of aggregates through the dryer. The dryer is equipped with flights or lifters on the interior for lifting and moving the wet aggregates through the hot gas stream as they pass through the cylinder.\nAs the gases pass through the dryer the aggregates become hot and dry thus liberating aggregate dusts that cling to the aggregates. Some aggregates liberate dust as they are tumbled through the rotary dryer. The dryer operates under a suction pressure as the dust particles are pulled through the dryer and into the coarse particle mechanical dust separation devices and finally into the fabric filters for dust collection. Those dust particles must pass from the rotary dryer to the dust collection devices generally using ductwork or large round pipes or rectangular ducts to move the moisture, products of combustion and dust to the dust collection devices. This industrial size external ductwork usually requires the use of a large lifting crane to assemble the ductwork from the rotary dryer to the dust collection devices and filter house needed to collect the hot dry aggregate particles. Moving the asphalt plant from one location to another becomes cumbersome as lifting equipment is needed to assemble and disassemble the plant rotary dryer from the dust collection equipment.\nMaking asphaltic mixes requires a recipe using a combination of sand and varying sizes of aggregates. Multiple bins are needed to contain the sand and aggregates. The bins are filled using construction bucket loading tractors from the site stored aggregates piles.\nThe aggregate bins are each fitted with an underlying variable speed belt feeder that proportions the needed flow of raw material to combine the hot mix recipe.\nThe raw aggregates are generally conveyed with a belt collection conveyor to a large vibrating screen to insure debris or oversize materials are eliminated from the flow of aggregates to insure a smooth road surface. A vibrating screen device of the size needed on the asphalt plant generally must be conveyed as a separate unit and mounted to a large wheeled trailer assembly in order for it to be conveyed for the portable plant. Bins of this size and capacity normally require a large individual wheeled trailer assembly as they are moved as a separate component of the portable asphalt plant. The movement of aggregates from the aggregate bins and oversize vibrating screen requires the use of a long portable conveyor belt assembly to move the aggregates to the rotary dryer for processing. These portable conveyors are fitted with truck axles and wheels as they are moved along with the portable asphalt plant.\nThe dust particle size separation and collection equipment needed on the portable plant is generally conveyed as a separate component with truck axles and wheels that also must be positioned on the plant site to mount the large connecting duct pieces between the rotary dryer and dust separation equipment. The large portable rotary dryers as well are conveyed as a separate component with truck axles. Ductwork connection and feed conveying plant components must be connected to the rotary dryer once it is in place at the portable plant site.\nOnce the hot mix asphalt is produced it must be conveyed away from the rotary dryer, stored and weighed and then transferred to a dump truck to be carried to the road paving site. Each of these devices must be made portable and carried along with the other plant components as a separate piece of equipment with truck axles as trailer assemblies.\nEach time a portable hot mix asphalt plant is moved, each site must be graded, leveled and compacted to insure the proper soil surface conditions exist to allow the equipment to be assembled and parked for alignment and safety. As described, the portable hot mix asphalt plant can easily contain ten individual plant components. Transporting and preparing the site represents time taken away from actual plant production; the cost to prepare the site, grading for each piece of equipment and the cost of additional site preparation materials such as soil and aggregates to prepare and level the site."}
{"text": "Desiccant cartridges containing desiccant particles are common in automotive air conditioning systems for dehydrating refrigerants. These desiccant cartridges are adapted to fit within canisters of integrated receiver/dryer (R/D) and accumulator assemblies.\nAlthough various forms of integrated R/D and accumulator assemblies have been suggested in the prior art, the types used in automotive air conditioning systems generally include an elongated canister having inlet and outlet ports communicating with the interior thereof. A desiccant cartridge including a desiccant cup for holding desiccant particles is positioned within the interior of the canister. A desiccant cap is secured to the cup to retain the particles therein.\nTypically, the inlet port is centrally located on a lid of the canister thereby allowing a center tube of the inserted desiccant cartridge to easily align and cooperate with the inlet port or a tube extending therefrom. This cooperation allows gas and/or fluid to flow into the canister and through the cartridge enabling the desiccant particles to remove moisture therefrom. More particularly, the gas or fluid enters the canister via the inlet port, travels through the tube extending therefrom, passes through the center tube, empties into the interior of the canister, flows through the desiccant cartridge, and finally exits the canister via the outlet port.\nIn contrast, some canisters are provided with an inlet port that is offset from the center such that a center tube desiccant cartridge must be accommodated in order to cooperate with the offset port or a side tube extending therefrom. One way to accommodate the combination including an offset inlet port and a desiccant cartridge having a center tube is to bend the side tube such that it connects the center tube of the cartridge with the offset port. However, this option is expensive and somewhat difficult to do.\nAlternatively, the desiccant cartridge with center tube can be replaced with a cartridge having a desiccant cup provided with an offset tube for aligning and cooperating with the offset inlet port, or side tube. However, one of the primary difficulties associated with assembling desiccant cartridges adapted to receive a side tube lies in orienting a hole, or aperture, in the desiccant cap to properly receive the offset tube of the desiccant cup. As such, assembly of current side tube receiving desiccant cartridges is inefficiently done by hand or with expensive automated assembly systems. Accordingly, desiccant cartridges provided with an offset tube typically are more expensive than the standard center tube cartridges.\nNotably, the desiccant cartridge of the present invention includes a desiccant cup having a center tube and a desiccant cap adapted to easily cooperate with both the center tube and an offset inlet port, or side tube, of a canister of an integrated receiver/dryer (R/D) or accumulator assembly thereby eliminating any need to properly orientate the cap on the cup and allowing for the use of simple, inexpensive parts and automated equipment for assembly thereof."}
{"text": "Eyelash coating compositions such as mascaras are generally makeup compositions, compositions to be applied over a makeup (also known as top coats), or cosmetic eyelash care compositions.\nMascaras are especially prepared according to two types of formulation: water-based mascaras known as cream mascaras, in the form of a dispersion of waxes in water; mascaras that are anhydrous or that have a low water content, known as waterproof mascaras, in the form of dispersions of waxes in organic solvents.\nThe present patent application more specifically concerns water-based mascaras.\nApplication of mascara makes it possible to increase the volume of the eyelashes and consequently to increase the intensity of the gaze. Many thickening or volumizing mascaras exist for doing this, the principle of which consists in depositing the maximum amount of material onto the eyelashes so as to obtain a volumizing (or charging) effect.\nIt is in particular by means of the amount of solid particles (especially waxes, which make it possible to structure the composition) that the application specificities desired for the compositions can be adjusted, for instance their fluidity or consistency, and also their thickening power (also known as the charging power or makeup power).\nThese solid particles are dispersed in the cream mascara by means of a surfactant system.\nAmong the standard emulsifiers or emulsifying systems, there are                cetyl phosphate, but the use of cetyl phosphate alone leads to aggregation of the pigments and also to coarse dispersion of the waxes; this frequently results in a “grey” rather than a black mascara;        emulsifying systems based on Steareth-20 and Steareth-2, but the use of these systems leads to very fluid mascaras whose consistency is unsatisfactory for a volumizing mascara;        emulsifying systems based on triethanolamine stearate.        \nOne problem addressed by the present patent application is a mascara in which not only the waxes but also the pigments are dispersed uniformly, the mascara having a texture that is thick enough to obtain a charging, volumizing deposit on the eyelashes, and having a satisfactory consistency allowing easy application to the eyelashes and smooth, uniform deposition."}
{"text": "Advancements in technology and the growing concern for environmentally efficient vehicles have led to the use of alternate fuel and power sources for vehicles. Electric vehicles or hybrid electric vehicles use energy storage systems (ESS) to provide power for various vehicle requirements. However, the vehicles must generate and recharge the ESS for continued usage.\nCommonly hybrid electric vehicles rely on torque and ESS state of charge to determine whether or not the generator needs to be activated to recharge the ESS. If the ESS state of charge falls below a set schedule then the generator is activated. It is desirable to create optimal charging schedules to recharge the ESS while providing minimal interruption to the operation of the vehicle. For an electric vehicle there is no way to recharge the ESS while the vehicle is drawing power. As a result, it is critical to be able to allocate the available energy to critical systems in an electric vehicle to ensure that the vehicle can reach the destination or next charging opportunity.\nTherefore, predicting the energy consumption required prior to and during the drive cycle in order to ensure energy availability and optimal charging schedules is desirable."}
{"text": "This section provides background information related to the present disclosure which is not necessarily prior art.\nThe Peltier effect is an effect in which a heat flux is created between the junction of two different types of materials. Thermoelectric modules (TEMs) are semiconductor devices which use the Peltier effect to transfer heat from one side of the TEM (the “cold side”) to the other side of the TEM (the “hot side”). The TEM, a solid-state Peltier device, effectively acts as a heat pump upon the application of a DC power source to the TEM. Heat is moved through the TEM, from one side to the other, with consumption of electrical energy, depending on the direction of the current. Such an instrument may also be called a Peltier cooler or heater, a thermoelectric heat pump, a Peltier heat pump, a solid state refrigerator, a thermoelectric cooler (TEC), or a thermoelectric module (TEM). TEMs can be used either for heating or for cooling (refrigeration), although in practice the main application is cooling.\nThermoelectric assemblies (TEAs) often include a pair of fan sinks that face outwardly in opposite directions and meet at one or more TEMs at their bases. A TEA may comprise a cold side fan sink, a hot side fan sink, and a plurality of TEMs disposed between the fan sinks. A fan sink typically includes a fan mounted on or adjacent to a heat sink, often through the use of a fan housing over the heat sink fins. The TEA may also include a foam gasket that surrounds the TEMs. The TEMs may be arranged in a single circuit or a series of parallel circuits. The fan sinks may both be in thermal contact with the TEMs heat sink bases."}
{"text": "1. Field of the Invention\nThe present invention relates to a signal distribution circuit for use in a converter for receiving cable television broadcast (hereafter referred to as a xe2x80x9cCATVxe2x80x9d converterxe2x80x9d).\n2. Description of the Prior Art\nConventionally, a CATV converter is used to permit a subscriber to view a desired program by selecting one among about 60 cable television broadcast signals (hereafter referred to as xe2x80x9cCATV signalsxe2x80x9d) that are transferred across a coaxial cable in a frequency range of, for example, 70 to 450 MHz (in the case of Japan). Such a CATV converter is installed on that end of the coaxial cable that leads to a television receiver of each subscriber.\nIn recent years, such a CATV converter has come to have enhanced functions that allow it to be connected to a videotape recorder for timer-recording, a second CATV converter for a second television receiver, an FM tuner, and the like. To achieve this, a modern CATV converter is additionally provided with output terminals for connection to a videotape recorder, a second CATV converter, an FM tuner, and the like, and also with a signal distribution circuit for distributing CATV signals among those output terminals.\nSpecifically, a CATV converter is configured as shown in FIG. 14. Reference numeral 31 represents a CATV signal input terminal at which the CATV converter receives the CATV signals transferred thereto across a cable. Reference numeral 32 represents a brancher (directional coupler) for branching the CATV signals received at the CATV signal input terminal 31 into signals that flow in a through direction and signals that flow in a branch direction. The signals output from the brancher 32 in the through direction are fed, on the one hand, to a CATV signal output terminal 40 for connection to an analog tuner, and, on the other hand, through a BPF (band-pass filter) 33 having a bandwidth of 70-75 MHz or 70-130 MHz that extracts, from those signals, downstream signals consisting of FSK (frequency-shift keying) signals having that band width, to a downstream signal output terminal 41 for connection to a tuner for receiving downstream signal.\nThe signals output from the brancher 32 in the branch direction are fed, through a variable attenuator 34 designed as, for example, an attenuator employing PIN diodes, then through two-stage wide-band amplifiers 35 and 36, to a distributor 37 that distributes the signals into two groups of signals. Here, the gain of the variable attenuator 34 is controlled by a gain control voltage that is applied to a control terminal (not shown) thereof. One of the two groups of signals output from the distributor 37 is fed to an extractor 38 composed of a band-pass filter that extracts VHF signals. The extracted VHF signals are, together with the UHF signals extracted by an extractor 39 composed of a band-pass filter from the signals fed in via a UHF signal input terminal 43, to a UHF/VHF signal output terminal 42 for connection to a UHF or VHF tuner of a television receiver. The other of the outputs from the distributor 37 is fed to an FM signal output terminal 44 for connection to an FM receiver. In this conventional example, a second CATV converter may also be connected to the output terminal 44.\nThe conventional example shown in FIG. 14 and described above is an example of a signal distribution circuit used in an analog STP (set-top box). In the case of a signal distribution circuit used in a digital STP, since it is connected to, as the succeeding stage, an analog/digital tuner and to a tuner for downstream signal reception, the signals are, before being fed to them, distributed into two groups of signals.\nIn a signal distribution circuit used in a conventional analog STB, much care is taken to transfer CATV signals with as small a loss as possible. To adapt this distribution circuit for digital-signal handling, it needs to be so configured as to permit an analog/digital tuner and a tuner for downstream signal reception to be connected thereto; specifically, an additional distributor needs to be provided after the band-pass filter 33 to distribute the signals into two groups of signals so that one of those outputs is fed to the downstream signal output terminal 41 and the other to a terminal (not shown) for connection to an analog/digital tuner. However, additionally providing this distributor causes a loss of 3 to 4 dB in the CATV signals, and thus degrades the C/N (carrier to noise) ratio. The distribution loss resulting from the insertion of this distributor can be compensated for by reducing the loss that is caused in the through direction by the brancher 32 provided in the preceding stage. However, reducing the loss caused in the through direction by this brancher 32 increases the loss in the branch direction, and thus degrades the NF (noise figure) of the signals that are extracted from the signals flowing in the branch direction so as to be fed to a television receiver and an FM receiver.\nMoreover, a signal distribution circuit used in a digital STB allows a subscriber to transmit upstream signals from his part to the CATV central station via the CATV signal input terminal 31. However, no conventional technique is known to be ready for this function. At this CATV signal input terminal 31, upstream signals need to have a level of 120 dBuV.\nAn object of the present invention is to provide a signal distribution circuit that has a simple circuit configuration but nevertheless achieves efficient distribution of received digital CATV signals.\nTo achieve the above object, according to one aspect of the present invention, a signal distribution circuit is provided with: a directional coupler for distributing input CATV signals into signals that flow in a through direction and signals that flow in a branch direction; a first extraction circuit for permitting, out of the signals output from the directional coupler in the through direction, only signals having predetermined frequencies to pass therethrough; a variable attenuator for adjusting the level of the signals output from the first extraction circuit to a predetermined level; a wide-band amplifier for amplifying the signals output from the variable attenuator; a first output terminal for connection to an analog/digital tuner for receiving CATV signals; a second output terminal for connection to a tuner for receiving downstream signals; a first distributor for distributing the signals output from the wide-band amplifier to the first and second output terminals; a filter for permitting, out of the signals output from the directional coupler in the branch direction, only VHF signals and FM signals to pass therethrough; an amplifier for amplifying the signals output from the filter; and a second distributor for distributing the signals output from the amplifier into two groups of signals.\nAccording to another aspect of the present invention, a signal distribution circuit is provided with: a CATV signal input terminal; a first extraction circuit for permitting, out of signals received at the CATV signal input terminal, only CATV signals to pass therethrough; a gain adjustment circuit for adjusting the level of the signals output from the first extraction circuit to a predetermined level; a wide-band amplifier for amplifying the signals output from the gain adjustment circuit; a first output terminal for connection to an analog/digital tuner for receiving CAIV signals; a second output terminal for connection to a tuner for receiving downstream signals; a first distributor for distributing the signals output from the wide-band amplifier into two groups of signals and feeding one of those two groups of signals to the first output terminal; a filter for permitting only VHF signals and FM signals to pass therethrough; a second distributor for distributing the other of the two groups of signals output from the first distributor to the second output terminal and the filter; a third output terminal; a fourth output terminal; and a third distributor for distributing the signals output from the filter to the third and fourth output terminals.\nAccording to another aspect of the present invention, a signal distribution circuit is provided with: a CATV signal input terminal; a first extraction circuit for permitting, out of signals received at the CATV signal input terminal, only CATV signals to pass therethrough; a first output terminal for connection to an analog/digital tuner for receiving CATV signals; a second output terminal for connection to a tuner for receiving downstream signals; a first distributor for distributing the signals output from the first extraction circuit into two groups of signals; a first circuit for adjusting the gain of one of the groups of signals output from the first distributor and feeding those signals to the second output terminal; a second circuit for adjusting the gain of the other of the groups of signals output from the first distributor; a filter for permitting only VHF signals and FM signals to pass therethrough; a second distributor for distributing the signals output from the second circuit into two groups of signals and feeding one of those two groups of signals to the first output terminal and the other to the filter; a third output terminal; a fourth output terminal; and a third distributor for distributing the signals output from the filter to the third and fourth output terminals."}
{"text": "Many computer systems, such as enterprise servers, support critical processes that should continue to operate despite faults and sub-system failures. Such systems cannot therefore be switched off when devices require removal or replacement. Moreover, removal of a device has the potential to cause damage to the system, unless it is carefully handled. The process of safely removing a device without switching off the system is known as online deletion.\nOnline deletion of I/O adapters has traditionally been implemented as a single step operation where the device driver attempts to disable the components in the adapter or render them inactive, an operation that is referred to herein as quiescing the adapter. If the quiesce operation fails, the online deletion operation fails. The quiesce operation may fail because one or more components in the adapter do not support the quiesce operation, because the operation returns a ‘failed’ status, or because the return status of the quiesce operation is unknown.\nIf the quiesce operation fails, then one or more components of the adapter may continue to initiate direct memory access (DMA) operations. Furthermore, if clean up/deallocation of datastructures or unmapping of DMA buffers is attempted as part of the online deletion operation, and one or more components of the adapter were unsuccessfully quiesced, then a subsequent DMA by the unsuccessfully quiesced components of the adapter will result in a ‘machine check’ condition, that generally prevents further operation of the system."}
{"text": "ISM (industrial, scientific and medical) radio bands are license-exempt bands, which were originally reserved internationally for the use of RF energy for industrial, scientific and medical purposes other than communications. However, in recent years these bands have also been shared with license-free error-tolerant communications applications such as Wireless LANs (WLAN) and cordless phones in the 915 MHz, 2.450 GHz, and 5.800 GHz bands. A main usage on these ISM bands is “WiFi.” While “WiFi” is not a technical term, the WiFi Alliance has generally enforced its use to describe only a narrow range of connectivity technologies including wireless local area network (WLAN) based on the IEEE 802.11, which is a set of standards carrying out WLAN communication in the 2.4, 3.6 and 5 GHz frequency bands: ISM band 2.4 GHz (WiFi 802.11b and 802.11g/n) and ISM band 5 GHz (WiFi 802.11a/n/ac).\nRecent growth in data traffic driven by mobile applications on smart phone devices, tablets, etc. has continued to strain the capacity of today's networks. Therefore, network operators are increasingly utilizing un-licensed WiFi spectrum to cope with such network congestion. Using unlicensed spectrum is a cost-effective mean to add needed capacity to today's networks and, given the limited availability and high cost of licensed spectrum, this trend is expected to accelerate further as traffic demand continues to grow.\nAccordingly, solutions for providing wireless connection switching that is as efficient and seamless as possible are needed."}
{"text": "Every year there are a significant number of people who suffer severe neck injuries from trauma. These injuries include cervical fractures, fracture dislocations combined with retropulsion of the disc and other major injuries. Also, each year many people undergo cervical spine surgery for degenerative diseases, especially degenerative spinal stenosis, which involves the removal of much of the bone which supports the cervical region of the spine resulting in instability of the cervical region.\nTreatment of these conditions has included traction either with a halter or with Crutchfield type tongs followed by application of a cast or brace. If surgery is necessary, the area of injury is often fixed with wire to allow fusion of the vertebrae in the affected region of the vertebral column. Often the treatment includes anterior decompression and fusion, or more recently, plates and screws have been used to immobilize the unstable region. Such plates may be used either anteriorly or posteriorly, or in a few cases, both anteriorly and posteriorly.\nWhile these fixation devices have been used effectively, they have a number of serious disadvantages. For example, the use of wires prevents flexion, as required for effective fixation, but does not prevent rotation nor extension of the cervical region, which leads to poor fixation of the vertebrae in this region. Currently, wires are nearly always used in combination with bone grafts. Sometimes the grafts are placed between the posterior elements of the vertebrae in such a fashion as to prevent extension, and since the wires prevent flexion, some stabilization is obtained. However, the stabilization is usually insufficient to avoid the necessity for the additional stability conferred by the use of halo devices, large casts or braces.\nIn the case where posterior plates and lateral bone mass screws are used, the stability is increased by simply screwing the plate into the bone. However, since the bone in the cervical area is relatively soft, the screws of such devices easily pull or cut out of the bone and, therefore, do not immobilize well. Even if there is no major failure of the device, such as the screws pulling out of the bone, the screws eventually work loose and plates, that initially prevented all motion, tend to loosen as the patient recovers from the surgery and becomes more active. As a result, halos, casts and braces are often used in conjunction with screw and plate fixation.\nNumerous fixation devices have been described, such as those described in U.S. Pat. Nos. 5,030,220 and 5,034,011, for use in the lumbar-sacrum regions of the spine. While these fixation devices have proven successful for use in the lumbar region, the physiology and structure of the cervical spine, C1 to C7, is very different from the lumbar-sacrum regions. For example, the bone screws, when placed in the lateral masses of the cervical vertebra, are usually placed at an angle from horizontal, whereas in the lumbar region they are placed in a horizontal plane. However, for proper stabilization and fixation, the hardware carried by the bone screws ought be such as to provide support and stabilization, generally along the axis of the cervical portion of the spine, rather than at an angle to it. It is desirable that a cervical fixation device be developed which would preferably take into account these considerations. In addition, it is desirable that a cervical fixation device be relatively compact and easy to install."}
{"text": "Numerous duct assemblies are known in the prior art for enclosing conduits. However, these duct assemblies have several disadvantages. For example, many of them require many additional pieces to join a plurality of duct assemblies together. In addition, many of these prior art devices do not provide any means for attaching another duct assembly thereto in both the X and Y directions, i.e., side by side or one on top of the other. Moreover, many of these prior art duct assemblies are very expensive to make because they require numerous parts.\nExamples of these prior art devices are disclosed in U.S. Pat. Nos. 3,693,664 to Schmunk; 3,727,644 to Kagan; 3,927,698 to Johannsen; 4,391,426 to Gothberg; 4,640,314 to Mock; 4,627,469 to Buard; and 4,720,768 to Schindele."}
{"text": "This application makes reference to, incorporates the same herein, and claims all benefits accruing under 35 U.S.C. xc2xa7119 from an application for Power Saving of a Portable Computer Using Human Sensing Device earlier filed in the Korean Industrial Property Office on Apr. 7, 1998 and there duly assigned Serial No. 12255/1998.\nThe present invention relates to portable computers, and more particularly to a portable computer having a power saving function that reduces supply power consumed in a display device.\nCompact and lightweight notebook and laptop computers have become increasingly popular due to their portability. These portable computers are typically powered by rechargeable batteries or alternatively by a power adapter which can be connected to an AC outlet. Often, it is desirable to operate the portable computers for lengthy periods of time in environments where AC power is unavailable.\nAs a system level solution, most portable computer systems usually adopt a power saving and power management system as provided in Advanced Power Management (APM) jointly developed by Intel (r) and Microsoft (r). This power management is to maximize available battery usage time. Basically, it automatically reduces the power consumption of parts of the computer currently not in use. The power management system provides various power saving function with the lapse of time set by the CMOS setup program, and manages the progress of power saving between full-on, doze, standby, and rest mode. The rest mode is the least power consumption level, which interrupts power supplies directed to microprocessor, display apparatus, hard disk drives, floppy disk drives, and other subsystems. Only supplied power is that with the subsystem which can restore the system to the full-on state.\nFurther, the power saving function may be performed by back-light timer with the lapse of time set by the CMOS setup program, which manages cut-off of power supply directed to a back-light unit of LCD display panel. When the system is not working for the preset time, the back-light unit of LCD display goes into power-off status. This power management is useful for portable computers having a LCD display panel to use less power while the system is running on battery power.\nHowever, in actual computer using condition, when the portable computer is powered on and a computer work is proceeding, the work may be interrupted by some matter and the user may leave the portable computer without turning off power or having the computer system go into the rest mode of the power management. In this case, the battery power is still consuming while the user is leaving the portable computer. Further, if the system power down times of the above power management have been set with long time intervals, the power saving function may be ineffective in saving the battery power. This requires modification of the power down times in the power management setup to suit work conditions whenever required. Also, the adjustment process is inconvenient for the user and conversely the modified power saving function may be unsuitable for a specific computer work.\nThus, there is a need in the portable computer art a more effective and simple power saving function that can minimize unnecessary power consumption during the computer use regardless of the power management setup condition.\nThe present invention is intended to solve the prior art problems, and it is therefore an object of the invention to provide a portable computer which can perform an instant power saving function in response to absence of the computer user during use of the computer.\nAccording to the present invention, the portable computer is provided with a sensing device which responds to absence of the computer user so that the computer system instantly goes into a power saving function that cuts off power directed to a back-light unit of the LCD display panel. A portable computer of the present invention comprises a sensor for sensing a computer user within the predetermined sensing range, said sensor being mounted in the housing of the flat panel display with facing the computer user, to produce a signal representative of absence of the user within the range; and a controller in response to said output signal of said sensor to produce a control signal which allows the supply power directed to the flat panel display to be interrupted.\nThe portable computer further comprises a basic input/output system (BIOS) which stores setup information regarding enabling/disabling the operation of said sensor. The sensor includes: a transmitter generating a sensing signal; and a receiver receiving the reflected sensing signal and responsive thereto outputting a signal representative of presence of the user.\nPreferably, the sensing signal consists of ultrasonic wave or infrared rays. Further, the controller responds to the output signal representating presence of the user, and determines whether a user password flag has been set in the system BIOS and if the password has been set a command is sent to the BIOS, resulting in generation of a message requesting input of the user password."}
{"text": "The present invention relates to a mobile terminal that can integrally move between an infrastructure network and an ad hoc network and to a method of controlling the same.\nAn infrastructure network in which networks such as internets are interconnected via a relay node such as a router as well as an ad hoc network which has no relay node and is a network temporarily configured of only terminals are well known as aspects of a network.\nIn order to establish communications by connecting a mobile terminal to a network, the mobile terminal must have (1) a terminal address used in a network to be connected, or a network address of the network itself and (2) an opposite communication party's address.\nWith respect to the item (1), the dynamic host configuration protocol, DHCP, (for example, refer to Douglas E. Comer, \"Internetworking with TCP/IP Volume I principles, protocols and architecture\", Third edition, 1995, Prentice-Hall, Inc.) is well known with an internet being an example of an infrastructure network. Moreover, with respect to the item (2), the domain name system, DNS, (for example, refer to Douglas E. Comer, \"Internetworking with TCP/IP Volume I principles, protocols and architecture\", Third edition, 1995, Prentice-Hall, Inc.) is well known. In the DHCP, a DHCP server is set up that holds a terminal address or infrastructure network address to be allocated to a terminal. Necessary values such as terminal or network addresses for connection are requested to the DHCP server when a terminal is connected to a network. The DHCP server provides parameters to the terminal without troubles due to duplicate terminal addresses, based on the request from the terminal. FIG. 15 depicts such an operation. When the terminal 1 and the terminal 2 are connected to a network, they respectively broadcast DHCP Requests containing desired information. FIG. 15 shows an example in which the terminals 1 and 2 request a terminal address. In response to the DHCP Request message, the DHCP server transmits a DHCP ACK message containing supply information to the request sources. In FIG. 15, the DHCP server transmits two DHCP ACK messages, one containing an address 1 and the other containing an address 2, to avoid overlapping of a terminal address at the terminals 1 and 2.\nIn the internet, the DNS (Domian Name System) is used to detect the correspondence between the terminal name of a terminal and the terminal address. The DNS introduces a DNS server that manages the correspondences between the terminal's names of terminals to be managed and the terminal addresses and predetermines the procedure of making inquiries from a terminal to the DNS server and the procedure of making inquiries between servers. In most cases, the address of the DNS server managing a terminal is set to the terminal itself. FIG. 16 depicts the case where the terminal 1 detects the terminal address of the terminal 2. The terminal 1 transmits a DNS Query message requesting the terminal address of the terminal 2 to the DNS server (the DNS server 1 in FIG. 16) registered. The DNS server 1 inquires the DNS server (the DNS server 2 in FIG. 16) managing the terminal address of the terminal 2. Then the DNS server 1 notifies the terminal 1 of the address 2 (DNS Reply message) when obtaining the terminal address of the terminal 2 (address 2).\nThe ad hoc network is configured of only terminals and does not have a server such as DHCP server or DNS server used in the infrastructure network. For this reason, even if the same procedure as that in the infrastructure network are used, communications cannot be accomplished by connecting the terminal to an ad hoc network. In order to accomplish communications according to the same procedure as that in the infrastructure network, a server function may be added to one of terminals connected to an ad hoc network. However, since this approach requires a special terminal with a server function, it is considerably poor in versatility. In order to improve the versatility, a procedure different from that in the infrastructure network may be defined in the ad hoc network so that the procedure can be selectively used according to an ad hoc network or infrastructure network to be connected. However, in the case of this procedure, the terminal must be reset when it moves between an ad hoc network and an infrastructure network according to the procedure for an infrastructure network and the procedure for an ad hoc network to be introduced thereto."}
{"text": "This invention relates to a method of determining the efficiency of an absorbent article having two absorbent layers that have a selected density, basis weight and each layer contains a superabsorbent. More specifically, this invention relates to a method of determining, without bench testing, the efficiency of a disposable absorbent article having two absorbent layers made from a wettable, fibrous material, each layer having a selected density, basis weight and a desired percentage of a superabsorbent.\nMany absorbent articles, especially disposable absorbent articles, like diapers, training pants, sanitary napkins, incontinence garments, and the like are constructed with two absorbent layers. Typically, the upper absorbent layer functions as a fluid intake layer and the lower absorbent layer functions as a fluid storage or reservoir layer. Normally, the superabsorbent material is present in the lower absorbent layer. This design allows the fluid which contacts the article to move downward through the upper absorbent and into the lower absorbent layer without being impeded. It is also common practice to construct the upper absorbent layer such that it has a lower density than the lower absorbent layer. This density gradient further facilitates downward fluid movement into the lower absorbent layer.\nNow, a method has been invented that can be used to determine the efficiency of an absorbent article having two absorbent layers before the absorbent article is actually manufactured. The absorbent article is constructed from a wettable, fibrous material such as cellulose fluff. The absorbent article has two absorbent layers each having a selected density, basis weight and a desired percentage of a superabsorbent in each of the two layers.\nBriefly, this invention relates to a method of determining the efficiency of an absorbent article having first and second absorbent layers each having a selected density, basis weight and a desired percentage of a superabsorbent. The method includes the steps of selecting a density for the first and second absorbent layers. A basis weight of at least 100 grams per square meter (gsm) is then selected for each of the first and second absorbent layers. The percentages of a superabsorbent present in each of the first and second absorbent layers are also selected such that the percentage of a superabsorbent in the first absorbent layer is greater than the percentage of a superabsorbent in the second absorbent layer. The selected densities and percentages of the superabsorbents are then inserted into the following formula:\nVxe2x89xa6[(100xe2x88x92% SAM2)xc3x97(d2xe2x88x92d1)]+[(% SAM1xe2x88x92% SAM2)xc3x97d2]\nwhere: % SAM1 is the percentage of a superabsorbent in the first absorbent layer;\n% SAM2 is the percentage of a superabsorbent in the second absorbent layer;\nd1 is the density of the first absorbent layer expressed in grams per cubic centimeters;\nd2 is the density of the second absorbent layer expressed in grams per cubic centimeters; and\nV is the desired value to form an efficient absorbent article.\nThe selected densities and the selected percentages of a superabsorbent in the first and second absorbent layers are then adjusted to obtain a value xe2x80x9cVxe2x80x9d that is at least equal to 2.\nThe general object of this invention is to provide a method of determining the efficiency, without bench testing, of an absorbent article having two absorbent layers each having a selected density, basis weight and each containing a desired percentage of a superabsorbent. More specifically, this invention relates to a method of determining the efficiency of a disposable absorbent article having two absorbent layers formed from a wettable, fibrous material, such as cellulose fluff and each absorbent layer having a selected density, basis weight and containing a desired percentage of a superabsorbent.\nAnother object of this invention is to provide a method of determining the makeup of a dual layer absorbent article using a mathematical formula.\nA further object of this invention is to provide a method of determining the percentages of a superabsorbent needed to be inserted into each absorbent layer of a two layer absorbent article.\nStill another object of this invention is to provide a method of economically designing a two layer absorbent article.\nStill further, an object of this invention is to provide an accurate and inexpensive method of designing a two layer absorbent article wherein each layer has a selected density, basis weight and a desired percentage of a superabsorbent.\nOther objects and advantages of the present invention will become more apparent to those skilled in the art in view of the following description and the accompanying drawings."}
{"text": "1. Field of the Invention\nThis invention relates generally to a single-sided modular cannula system for delivering oxygen and/or filtered air. More particularly, the invention relates to a single-sided modular portable system for delivering oxygen and/or filtered air to persons in a comfortable and convenient manner while providing additional options for use in combination with telecommunications and audio devices.\n2. Description of Related Art including information disclosed under 37 CFR 1.97 and 1.98\nClean air is an important part of maintaining health in an otherwise healthy person. Environments that contain air-borne pollutants and infectious agents have received much publicity and awareness. Second hand cigarette smoke, carbon monoxide, SARS, dust, pollen and car exhaust, among other such pollutants can cause respiratory discomfort, damage or inefficiency.\nSome of these effects can be temporary, as when exposed to an allergen where the effect disappears when the allergen is removed, and some of these effects can be permanent, as when second-hand cigarette smoke causes cancer. It is beneficial to an otherwise healthy person to maintain a high quality supply of breathable air by either filtering the ambient air before inhalation, or providing an alternate source for the breathable air that excludes a large fraction of the ambient air.\nThis effect is well known for unhealthy people, such as an emphysema patient who pulls a cart with an oxygen tank that supplies oxygen to an area around the nose or mouth to enhance the oxygen uptake of damaged lungs. Another example is the oxygen masks used in hospitals to provide a similar effect for the same or other medical purpose. These systems provide higher quality air in a manner that is not easily portable for an individual, and certainly not designed to be fashionable or attractive.\nA personal oxygen and/or filtered air delivery system must further meet certain functional requirements. It must be able to remove some potentially, perceived or actually harmful fraction of the ambient air, or must be able to provide at least a portion of the inhaled air by displacing at least a portion, if not all, of the ambient air with either air filtered by an air filtration pump, or air from a compressed gas tank of appropriate dimension so as to be worn on the body. In the alternative it may be required to do both functions, with some portion of the inhaled air coming through the filter, some portion coming from the gas tank, and the remaining air, if any, coming from the ambient air.\nFor the purposes of this application, higher quality air refers to air with either more desirable characteristics, air which has fewer undesirable characteristics or air which has both more desirable or fewer undesirable characteristics. Percent O2 content, pollen, particles, aromatic compounds, gaseous and particulate carbon compounds including hydrocarbons, nitrogen containing compounds, carbon monoxide, ozone, viral infectious agents, bacterial infectious agents, sulfur containing compounds, dust, soot, smoke, smog, and many other compounds can, among many other factors, be factors used when determining characteristics of air quality. Desirable and undesirable characteristics can be left to the objective determination of government and industry agencies or subjective determination of each individual user of this system.\nAs air quality in some cities declines, and as the perception that the number and quantity of harmful components in the air is increasing, there is a growing group of people who wish to breath higher quality air but are unwilling to pull a tank of higher quality air around in a cart, or unwilling to wear a standard mask or nose-piece normally associated with an unhealthy person, and thus convey a negative body image to others or to themselves.\nFurther, there has developed a desire by some people to stop periodically in commercial locations called oxygen bars, or commercial locations, like health clubs, that offer oxygen bars in addition to other services. These commercial locations provide customers with, among other services, an opportunity to temporarily inhale breathable air that can contain higher concentrations of oxygen.\nSuch breathable air may also be enhanced with aromas pleasing to the customer. Currently these customers use a small flexible hose called a cannula to deliver the oxygen enriched breathable air into their nostrils. As part of the experience of being in an oxygen bar and social setting, it is desirable to portray the experience as glamorous or fashionable in addition to being functional.\nAccordingly, it is desirable for the cannula to be incorporated into other devices normally appearing on or around the face of a person to enhance the appearance of the medical oxygen user, or to discretely disguise the functional aspect of the cannula.\nVarious single-sided modular embodiments of the present invention are directed to addressing various needs in connection with ensuring that higher quality air and/or oxygen is delivered to the nose and or mouth area of a person wearing the system disclosed in this invention. These various embodiments encourage those who need oxygen to be as active and socially engaging as they were previous to being prescribed medical oxygen. These single-sided modular options attaches the oxygen and/or filtered air delivery system via proprietary adaptors to either an ear-piece, aviators telecommunications headset, helmet, cap, visor or hats, as well as to the frame of eye glasses.\nAs people feel the need to maintain communications in a near constant manner, they often make use of portable wireless earpieces. In a similar manner, many people choose to take advantage of compact headphones for use with personal music systems. The combination of these devices with oxygen/air breathing devices represents new flexibility and freedom for those who desire to take advantage of both systems in a comfortable and convenient manner.\nVarious patents address systems that are designed to permit the wearer to breathe oxygen or purified air. Sanders, U.S. Pat. No. 5,582,164, disclose a portable gas delivery system that includes a gas storage assembly that is connected to a nose piece using flexible tubing. The gas storage assembly includes a strap that enables it to be readily carried by the person using this system.\nAmRhein, U.S. Pat. No. 4,996,983, discusses a portable oxygen delivery system in which gas storage containers are incorporated into the temples on eyeglasses. Air delivery tubes extend around the wearer's ears to a nose piece, prongs or portal that is placed in the wearer's nose.\nPeppler, U.S. Pat. No. 5,193,534; and Timmons et al., U.S. Pat. No. 4,559,941, both describe incorporating an oxygen delivery system into eyeglasses. The oxygen is delivered from a storage container to the temples on the eyeglasses. The oxygen is then routed to a nosepiece through the eyeglasses.\nKoch et al., U.S. Pat. No. 4,465,067, discloses a system for delivering oxygen. This system has a configuration that is similar to eyeglasses except that it does not include any lenses. Oxygen is delivered to the temple portion that wraps around the wearer's ear and then to a nose piece that is positioned proximate the wearer's nose.\nIzuchukwu et al., U.S. Pat. No. 6,526,968, discusses an air delivery system that is incorporated into a utility belt that is worn by the user. The utility belt includes a pack for storing the compressed gas and a port for connecting to a mask using flexible tubing.\nIzuchukwu et al., U.S. Pat. No. 6,510,859, discloses an emergency breathing apparatus that includes a hood, which is placed over the user's entire head. Air is delivered to the hood from a storage pack. While such a system is acceptable for emergency situations, the fact that it covers the user's entire head limits the ability for the user to perform many activities while wearing this device.\nMcDonald, et al., U.S. Pat. No. 6,595,207, describes an oxygen diffuser for a lightweight oxygen delivery system for a patient, the oxygen delivery system being of the type comprising a mount for seated engagement on a patient's head or ear, an elongated tubular boom for oxygen delivery secured at one end to the mount and having the diffuser secured to the other end, the diffuser to deliver oxygen passed through the boom to a space in the vicinity of the patient's nose and mouth, the diffuser comprising a body having a wall, the interior surface of which wall is of generally concave configuration, circumscribing a centrally positioned oxygen outlet so as to direct the flow of oxygen from the outlet generally towards the patient's nose and mouth; and a baffle seated over the oxygen outlet so as to assist in mixing of oxygen with ambient air and avoid a direct flow of oxygen towards the patient's face.\nDemers, et al., U.S. Pat. No. 6,619,288, describe a breathing mask for delivering oxygen to a patient. The breathing mask has a headset for seating about the cranial region of the head of the patient, with a hollow gas delivery arm coupled to the headset. The mask also has a nosepiece that is coupled only to the hollow gas delivery arm and that delivers oxygen to the nostrils of the patient. Finally, the mask has a source of oxygen for coupling oxygen to the hollow gas delivery arm.\nMcCombs, et al., U.S. Pat. No. 6,065,473, describes a non-contact gas dispenser comprising a head set apparatus, a gas source and a conduit for directing the desired gas to a region proximate to the user's nose and mouth. The gas source preferably is a pressure swing adsorption apparatus that allows the user to select one of at least two pre-determined settings of product gas, each setting having a distinct concentration and flow rate different from the other setting(s).\nMartin, James F. et at., U.S. Patent Application No. 2006/0042637, describes an oxygen delivery device for regulating the flow of oxygen from an oxygen source to a respiratory cannula located on the face of a patient. A host controller coupled to the oxygen delivery device accepts an input from a medical monitor relating to the blood oxygen saturation percentage of the patient. The host controller then regulates the rate of oxygen supplied to the patient based in part upon the blood oxygen saturation percentage of the patient. A high rate of oxygen is supplied to the patient when the controller detects a blood oxygen saturation level below a predetermined percentage and a low rate of oxygen is supplied when the controller detects a blood oxygen saturation level above a predetermined percentage in an order to minimize wasted oxygen."}
{"text": "The present invention relates generally to apparatus for folding a continuous length of web material along a longitudinal fold line, and more particularly to a web folding apparatus having novel features which enable highly efficient folding of a continuous length web material without forming wrinkles or gussets and without marking the web, and which facilitate on-going adjustment to accommodate different type web materials.\nApparatus for folding continuous length web material along a longitudinal fold line as the web material is moved longitudinally are generally known. Such apparatus find particular application in the folding of paper web material as it is passed from a printer to a cutter and sheeter. One type of apparatus for folding web material employs a frame on which transverse entrance and exit rollers are mounted for passage of a web material from the entrance to the exit roller and between which a fold forming mechanism is supported for cooperation with the entrance and exit rollers to fold the web material longitudinally. A significant problem which exists with known web folding apparatus is the formation of wrinkles or gussets at the fold line and marking of the web material in the area of the fold. While the problem of web marking is of particular significance in web folding apparatus employed in the printing industry, the problem of wrinkling or gusseting adversely affects the folding of substantially any web material which is folded generally longitudinally thereof.\nAnother significant problem found in the known prior web folding apparatus is that they are generally incapable of adjustment during movement of the web material, particular as affects adjustment of the fold former plow, so that shutdown of the apparatus is required for making fine adjustments as may be required to reduce problems of wrinkling or gusseting and marking of the web material, thus adding significantly to the cost of operation.\nStill another problem which must be confronted in the longitudinal folding of many web materials is the tendency for the folded portion of the web to unfold following initial folding. While prior known web folding apparatus have attempted to address this problem, the various means suggested, such as goose neck type arm supported rollers and the like for engaging the folded web generally at the fold line, have not been found satisfactory."}
{"text": "Digital video compression ranges from coding still images to moving video for photographic, broadcasting, streaming, and conferencing applications. In traditional transform coding, pictures consisting of pixels are partitioned into blocks; each block of pixels is transformed; the transform coefficients are ordered and quantized; and some form of variable length coding (VLC, including arithmetic coding) is used to represent the series of quantized transform coefficients of each block. In many current standards, the blocks are of a fixed size, e.g., 8×8 for JPEG, MPEG1, MPEG2/H.262, MPEG4, H.261, and H.263. While some coding methods, e.g., the high profile of H.264/AVC, allow for two block sizes (4×4 and optionally 8×8) within the same picture, the VLC method for all blocks is the same for the one picture. In all of these standard coding methods, the VLC method is either two-dimensional variable length coding (2DVLC) or one of its many variations, such as Context-based Adaptive Binary Arithmetic Coding (CABAC), Context Adaptive Variable Length Coding (CAVLC) and Context-Based 2D Variable Length Coding (C2DVLC)."}
{"text": "This disclosure relates generally to virtualization in computing systems, and more particularly to synchronization, replication, and migration of micro-partitions. A micro-partition in a high availability solution may regularly save its processing state and its memory state to an alternate host machine to provide failover support in the event of a primary host machine failure, and may require the availability of processing resources associated with the alternate host machine in addition to resources associated with the primary host machine."}
{"text": "A wide variety of fibrous materials have been employed in tobacco smoke filter elements. However, the choice of materials for use in production of such filters has been limited because of the need to balance various commercial requirements. A very important property of a tobacco smoke filter is obviously its filtration efficiency, i.e., its ability to remove selected constituents from the tobacco smoke. However, the range of filtration efficiency has had to be compromised in order to satisfy other commercially important factors such as resistance to draw, hardness, impact on taste, and manufacturing costs.\nCellulose acetate has long been considered the material of choice in the production of tobacco smoke filters, primarily because of its ability to provide commercially acceptable filtration efficiency, on the order of about 50%, without significantly detracting from the tobacco taste, low resistance to draw, and filter hardness desired by the majority of smokers.\nA significant component of the commercially desirable \"taste\" is provided by the standard plasticizers utilized in the production of filter elements from cellulose acetate fibers, usually triethylene glycol acetate or glycerol triacetate (\"triacetin\"). In conventional cigarette filter manufacturing, the plasticizer is commonly applied to the cellulose acetate fiber by spraying or wicking using art-recognized techniques. The tendency of the plasticizer to migrate toward the center of conventional cellulose acetate fibers reduces the level of plasticizer at the fiber surface, minimizing its taste-enhancing capability and limiting the shelf life of plasticized tow fibers before being processed into filter rods. The plasticizer is therefore usually added to the tow during the manufacture of the filter rods.\nCellulose acetate fiber plasticized in this manner and wrapped with paper into rod-like forms become bondable at the fiber contact points, enabling the formation of relative self-sustaining, elongated filter rods in two to four hours. This process can be accelerated by the application of gases at elevated temperatures simultaneously with the formation of the filter rod. Filter rods produced in this manner provide a tortuous path for the passage of tobacco smoke when discrete lengths of such material are utilized as tobacco smoke filter elements.\nFiltration efficiency can be increased significantly through the use of small fibers which provide increased fiber surface area at the same weight of fiber. Solvent spun cellulose acetate fiber is commercially available only in fiber sizes down to 13 microns in diameter. To obtain finer cellulose acetate fiber, e.g., 10 microns or less, melt spinning of plasticized cellulose acetate resin would be required; however, the level of plasticizer necessary to directly spin such fine cellulose acetate fibers would render the resultant fibers very weak and commercially useless. Melt spun cellulose acetate of a larger diameter, which would require less plasticizer, would have to be drawn and crimped to produce such fine fibers for use in tobacco smoke filters. Unfortunately, melt spun cellulose acetate fibers can only be commercially drawn at relatively low draw ratios before the fibers break during processing. The inability to form and process very fine fibers of cellulose acetate places practical limits on the filtration efficiency capabilities of this material in the production of tobacco smoke filters.\nFurther, and very important commercially, by comparison with other polymeric materials such as the polyolefins, cellulose acetate is relatively expensive, costing, for example, on the order of more than three times as much as commercially available polypropylene in resin form. While attempts have been made to utilize other less expensive and more easily processed polymeric materials such as polypropylene in lieu of cellulose acetate in the manufacture of tobacco smoke filters, such efforts have been almost universally abandoned on a commercial level, primarily because of the undesirable impact of such materials on the taste properties of tobacco smoke. Also, such use is generally limited by the inability to easily bond the fibers in order to obtain the desired filter hardness at required resistance to draw.\nAnother problem with commercially available tobacco smoke filters, particularly cigarette filters, currently on the market is the difficulty in disposing of such materials after use. By bonding highly crimped cellulose acetate fibers at their contact points, conventional cigarette filters are designed to provide a significant volume of interstitial space for the passage of smoke. The bonded contact points of such filter elements degrade very slowly under normal environmental conditions resulting in high volume, long life, environmentally undesirable litter."}
{"text": "With the advent of the internet and online electronic commerce (e-commerce), financial transaction programs that allow users to seamlessly purchase products, transfer and receive funds over an internet connection have been in high demand.\nTraditional methods of executing financial transactions have been limited to a user providing his or her credit card, debit card, or checking account number on a commercial website, or using checks, money orders and other forms of paper-based payments. However, these means of executing financial transactions are often cumbersome, slow, and inconvenient, requiring a user to remember a multitude of account numbers, login data and passwords. This often results in significant time delays for payment processing. Furthermore, security and fraud concerns are prevalent. For instance, a user is often reluctant to provide sensitive credit card or debit card information over an internet connection, regardless of how “secure” an internet connection claims to be.\nRecent financial transaction programs have emerged as a means for a user to pay for purchases, transfer money, receive money (if the user is a merchant), store shipping addresses, and set up multiple financial accounts (e.g. checking or savings, credit card, debit card) all with one single login and password. Security and fraud concerns are also mitigated by means of online financial security precautions, encryption methods, and anti-phising programs that are inherent in online, internet-based financial transaction systems.\nMerchants may decide to use such a financial transaction program in order to receive funds from customers who purchase products from their commercial websites. There are several “payment mechanisms” that enable merchants to receive funds from customers. For instance, certain mechanisms enable customers to enter their credit card, debit card or go to an external login site to enter their account information, pay for the product, and then proceed to checkout.\nHowever, it is difficult for merchants to integrate these payment mechanisms without significant experience with the associated application programming interface (API) or knowledge of computer code. Merchants or vendors have to understand the API, have to know how to code in the specific provided language (i.e. ASP, ASP.NET, JSP, Java, PHP) and have to know how to develop complicated code that interfaces with a plurality of servers and databases during an online financial transaction between a customer and a merchant.\nCurrently, there are no user-friendly programs that present a simple contextual page flow and that also simultaneously generate the necessary code to integrate into a merchant's commercial website.\nTherefore, there is a need in the art for a simple, user-friendly and customizable method and system that allows a merchant to seamlessly integrate payment mechanisms inherent to financial transaction programs into their commercial websites, without having to know an API library or the details of a particular code (i.e. ASP, JSP, PHP, .Net, Java)."}
{"text": "Many chemotherapeutic methods are now available to be used in the treatment of cancer. One of the most successful methods is the use of anti-mitotic agents which interfere with the assembly or disassembly of microtubules. Since microtubule assemble and disassemble is necessary for mitosis, inhibition of either the assembly or disassembly of microtubules interferes with cell proliferation. Thus, compounds that inhibit the assembly of microtubule are useful in treating diseases or conditions which are caused or exasperated by rapid or abnormal cell proliferation, such as cancer.\nSeveral anti-mitotic agents have had considerable clinical success. For example, the following vinca alkaloids which inhibit microtubule assembly have proved clinically successful: Vincristine has been successfully used to treat hematological malignancies and non-small-cell lung carcinoma; Vinblastine has been successfully used to treat hematological malignancies, testicular carcinomas and non-small-cell lung carcinoma; and Vinorelbine has been successfully used to treat hematological malignancies, breast carcinomas and non-small-cell lung carcinoma. In addition, taxanes which inhibit microtubule disassemble have also proved to be clinically successful. For example, Paclitaxel has been successful in treating breast, ovarianand non-small-cell lung carcinomas; and Docetaxel has been successful in treating breast and non-small-cell lung carcinomas.\nDespite these successes, available anti-mitotic agents are inadequate for a number of reasons. For example, paclitaxel, docetaxel and vincristine are associated with significant neuropathy which can limit their use in repeat courses of therapy. In addition, both the vinca alkaloids and taxanes are good substrates for the 170 kDa P-glycoprotein (Pgp) efflux pump found in most multi-drug resistant cells. This protein pumps a drug out of the tumor cells causing the tumor cells to become resistant to treatment. Once a patient's cancer has become multi-drug resistant, there is typically little that can be done to halt or retard further progression of the disease.\nThere is therefore still a need for new drugs which overcome one or more of the aforementioned shortcomings of drugs currently used in the treatment of cancer. Desirable properties of new anti-cancer drugs therefore include a good therapeutic index, efficacy against tumors that are currently untreatable or poorly treatable, efficacy against multi-drug resistant tumors and/or reduced side effects."}
{"text": "1. Field of the Invention\nThe present invention relates to a demodulation circuit and a demodulating method for demodulating a digital-modulated signal.\n2. Description of the Related Art\nFIG. 1 is a block diagram of a common QPSK demodulation circuit. A QPSK-modulated signal is received by a tuner not shown in the attached drawings, and converted into an intermediate frequency signal IFin. The intermediate frequency signal IFin is converted by an A/D converter 12 into digital data at the timing synchronous with the sampling clock with a predetermined frequency, and further converted by an IQ separator 13 into a baseband signal formed by an I signal and a Q signal.\nThe IQ separator 13 includes multipliers 13a and 13b for multiplying the output of the A/D converter 12 by a sine wave and a cosine wave having the frequencies equal to the center frequency of the intermediate frequency signal IFin. Low pass filters (LPF) 14a and 14b remove the upper frequency components of the output of the multipliers 13a and 13b, and the output is transmitted to interpolaters 15a and 15b. \nEach of the interpolaters 15a and 15b is constituted by an FIR (finite impulse response) filter, a thinning circuit, etc., and the data of the original sampling point (symbol timing) is obtained by interpolation from the received data sampled by the A/D converter 12.\nEach of the Root Nyquist filters 16a and 16b is constituted by a low pass filter, and restricts the band of the output signal of the interpolaters 15a and 15b. \nA rotor 17 is constituted by a butterfly circuit, etc., advances or delays the phase of a regenerated carrier wave according to the phase difference signal output from a carrier recovery circuit 18, and thereby allows the regenerated carrier wave to be synchronous with the carrier wave of a received signal.\nThe carrier recovery circuit 18 calculates the phase difference between the signal point on the I and Q phase plane obtained from the I signal and the Q signal and the normal signal point, and outputs the signal in the phase difference decreasing direction to the rotor 17.\nA timing recovery circuit 19 calculates the amount of shift of the sampling timing, and output an interpolating tap coefficient to the interpolaters 15a and 15b. \nFIG. 2 shows the configuration of a timing recovery loop. The timing recovery loop is constituted by the interpolaters 15a and 15b and the timing recovery circuit 19.\nThe timing recovery circuit 19 includes a phase comparator 21, a loop filter 22, a numerical controlled oscillator 23, a thinning control unit 24, and a tap coefficient arithmetic unit 25. The interpolaters 15a and 15b includes Fir filters 26a and 26b, and thinning circuits, 27a and 27b. \nThe phase comparator 21 determines whether the value of a received signal at each sampling timing indicates a monotonous increase, a monotonous decrease, or none of them, and outputs a signal depending on a discrimination result.\nThe loop filter 22 is a completely secondary loop filter, and includes a multiplier 28 constituting a low pass filter, a multiplier 29 constituting a high pass filter, an adder 30 for adding the output of the multiplier 28 to the output of the integrator 31, a limiter 32, an integrator 31, and an adder 33 for adding the output of the integrator 31 to the output of the multiplier 29. The integrator 31 is constituted by a flip-flop, etc. The coefficient α provided for the multiplier 28 is a gain adjustment coefficient for a low pass filter, and the coefficient β provided for the multiplier 29 is a gain adjustment coefficient for a high pass filter.\nThe numerical controlled oscillator (NCO) 23 includes a flip-flop 34 and an adder 35 for adding the output of the flip-flop 34 to the output of the loop filter 22. The numerical controlled oscillator 23 outputs digital oscillation frequency data according to a signal indicating the shift of the sampling timing output from the loop filter 22.\nThe tap coefficient arithmetic unit 25 provides a tap coefficient for advance or delay of a phase for the Fir filters 26a and 26b depending on the oscillation frequency data output from the numerical controlled oscillator 23, and outputs the I signal and the Q signal at the original sampling point from the Fir filters 26a and 26b. \nThe thinning control unit 24 controls the thinning circuits 27a and 27b, and thins data at an unnecessary sampling point in the data at each sampling point output from the Fir filters 26a and 26b. \nNext, the operation of the phase comparator 21 of the above-mentioned timing recovery loop is explained by referring to the operation explanation shown in FIG. 3.\nFIG. 3 shows the analog waveform of the input of the phase comparator 21. FIG. 3A shows the case in which the value of the input signal of the phase comparator 21 at each of the sampling time t−1, t and t+1 indicates a monotonous increase, and FIG. 3B shows the case of a monotonous decrease.\nIt is determined whether the values of the input signal indicate either a monotonous increase, based on the values of the input signal of the phase comparator 21 at each of the sampling time t−1, t and t+1, which are d(t−1), d(t), and d (t+1), a monotonous decrease, or any other cases, a predetermined arithmetic operation is performed on the value of an input signal at each sampling time based on the determination result, and outputs the arithmetic result as a phase determination result.\nThe phase comparator 21 outputs a value obtained by the following equation when the values of the adjacent input signal at the sampling time t−1, t, t+1 indicate a monotonous increase.−2×(d(t)+{d(t+1)−d(t−1)}/2)\nThe phase comparator 21 outputs a value obtained by the following equation when the values of the input signal at the sampling time t−1, t, t+1 indicate a monotonous decrease.2×(d(t)+{d(t+1)−d(t−1)}/2)\nFurthermore, the phase comparator 21 outputs “0” when the values of the input signal at the sampling time t−1, t, t+1 do not indicate a monotonous increase nor a monotonous decrease.\nFIG. 4A shows the output of the loop filter 22 when the frequency error between the original sampling point (hereinafter referred to as a symbol timing) and the sampling clock is small. FIG. 4B shows the output of the loop filter 22 when the frequency error against the symbol timing is large.\nWhen the frequency error is small, and the values of the input signal indicate a monotonous increase or a monotonous decrease, the phase comparator 21 outputs a value indicated in the above-mentioned equations, and the integrator 31 of the loop filter 22 integrates the values. As a result, for example, as shown in FIG. 4A, the output of the loop filter 22 increases until it reaches the convergence point of the timing recovery loop, and the capturing operation is completed at the convergence point.\nOn the other hand, when the frequency error is large, the signal points are distributed at random on the I and Q phase plane. Therefore, the average value of the output of the phase comparator 21 is approximately 0. As a result, the output of the loop filter 22 is fixed to a specific value, for example, as shown in FIG. 4B, close to “0”, or is fixed to a specific value.\nFIG. 5 shows the value of the limiter 32. When the average value of the output of the phase comparator 21 is approximately 0, the value of the limiter 32 to which an integration result is input is also a constant value.\nThe characteristic shown by the dotted line in FIG. 4B indicates the desired output characteristic of the loop filter 22 for reaching the convergence point, but the above-mentioned timing recovery circuit 19 cannot realize the capture characteristic as such.\nThat is, when the timing error is large, the output of the loop filter 22 is fixed to a specific value, and it is hard for the timing recovery circuit 19 to maintain the synchronization of clock timing.\nFIG. 6 shows the configuration of the carrier recovery loop for synchronization of the regenerated carrier wave with the carrier wave of a modulated signal.\nThe carrier recovery loop includes the rotor 17 and the carrier recovery circuit 18. The carrier recovery circuit 18 includes a phase comparator 41, a loop filter 42, a numerical controlled oscillator 43, and a Sin/Cos table 44.\nThe phase comparator 41 determines the advance or delay of the phase of the signal point on the I and Q phase plane obtained from the I signal and the Q signal to the phase of the normal signal point, and outputs a phase error signal for amendment of the advance or delay of the phase.\nThe loop filter 42 is a completely secondary loop filter, and includes a multiplier 45 constituting a low pass filter, a multiplier 46 constituting a high pass filter, an adder 47 for adding the output of the multiplier 45 to the output of an integrator 48, a limiter 49, the integrator 48, and an adder 50 for adding the output of the integrator 48 to the output of the multiplier 46. The integrator 48 is constituted by a flip-flop, etc. The coefficient α provided for the multiplier 45 is a gain adjustment coefficient for a low pass filter, and the coefficient β is a gain adjustment coefficient for the high pass filter.\nThe numerical controlled oscillator 43 includes a flip-flop 51 and an adder 52 for adding the output of the flip-flop 51 to the output of the loop filter 42. The numerical controlled oscillator 43 outputs oscillation frequency data for amendment of the frequency error (phase error) of a carrier.\nThe Sin/Cos table 44 is a table for generation of a sine wave and a cosine wave corresponding to the oscillation frequency data output from the numerical controlled oscillator 43.\nThe rotor 17 includes multipliers 53 and 54 for multiplying an I signal by the sine wave and the cosine wave output from the Sin/Cos table 44, multipliers 55 and 56 for multiplying a Q signal by the sine wave and the cosine wave, an adder 57 for adding the output of the multiplier 53 to the output of the multiplier 56, and an adder 58 for adding the output of the multiplier 54 to the output of the multiplier 55. (The output of the multiplier 54 is multiplied by “−1” and input to the adder 58.)\nThe rotor 17 multiplies the I signal and the Q signal by the sine wave and the cosine wave output from the Sin/Cos table 44 to rotate the signal point by the primary transform, and outputs the I signal and the Q signal at a desired sampling point.\nFIG. 7 shows a signal point of the I and Q phase plane (constellation point). The point indicated by the black dot shown in FIG. 7 is a normal signal point, the point indicated by a white dot is a signal point indicating a minus phase difference, and the point indicated by a dot with diagonal lines is a signal point having a plus phase difference. The counterclockwise rotation is a positive rotation and the clockwise rotation is a negative rotation.\nThe phase comparator 41 outputs a minus error signal when the position of the signal point on the I and Q phase plane determined according to the I signal and the Q signal output from the rotor 17 is detected in a predetermined range of the plus side relative to the position of the normal signal point, and outputs a plus error signal when it is detected in a predetermined range on the plus side on the I and Q phase plane. The phase comparator 41 outputs “0” when the position of the signal point depending on the I signal and the Q signal is detected in the position of the normal signal point.\nIn the carrier recovery loop, when the frequency error of a carrier is large, a signal point is distributed at random on the I and Q phase plane. Therefore, the average value of the output of the phase comparator 41 is approximately 0. Therefore, there has been the problem that, like the above-mentioned timing recovery loop, the output of the loop filter 42 is fixed to a specific value, and no carrier can be captured.\nThe patent document 1 describes providing a limiter for outputting an initialization value at the minimum level when the input value becomes larger than the limiter value in the feedback route to prevent pseudo synchronization.\nThe patent document 2 describes integrating a phase error signal, detecting synchronization by comparing an integrated value with a threshold, and stopping frequency sweep when synchronization is detected.\n[Patent Document 1] Japanese Patent No. 2885058\n[Patent Document 2] Japanese Patent Publication No. H5-30098"}
{"text": "1. Field of the Invention\nThe present invention relates to a vehicle body structure. More particularly, the present invention relates to a vehicle body structure configured to ensure both the strength of the vehicle body and the protection of the passengers.\n2. Background Information\nOne example of a conventional vehicle body structure is disclosed in Japanese Laid-Open Patent Publication No. 9-202268 (page 2 and FIG. 1) that uses extruded light alloy materials for main framework members of the vehicle body to reduce the weight of the vehicle body and improve the protection of the passengers.\nIn view of the above, it will be apparent to those skilled in the art from this disclosure that there exists a need for an improved vehicle body structure that ensure both the strength of the vehicle body and the protection of the passengers. This invention addresses this need in the art as well as other needs, which will become apparent to those skilled in the art from this disclosure."}
{"text": "The invention relates generally to sound absorbent material. In a particular embodiment, it relates to a thin, porous metal sheet which, when used with an air space, functions as a sound absorbent material for industrial applications.\nRecently, increased attention has been focused on noise pollution, including that generated by the use of machinery, such as drills, lathes and the like. Noise is a problem to the operator of the machine, as well as to those in the area where the machine is being operated. As a result, efforts have and are being made to prevent or to reduce the noise level of machinery in order to provide a safer, quieter work area.\nAttempts have been made to silence or to reduce machine noise by lining the inside of the machine, or housing surrounding the machine, with a material which will soak up or reduce the noise. It is important that the sound absorbent material be relatively inexpensive in addition to being an efficient and effective sound absorber. If the material is too expensive, the expense will make its use prohibitive even though it is efficient and effective, especially considering the number of machines which would use the material. Thus, an inexpensive material would find greater use, even where it was inefficient or would eventually become ineffective and have to be replaced, because of its low cost.\nA typical lining material for machinery is open-cell polyurethane foam, which is a resistive sound absorber. It is a relatively inexpensive material, and it will reduce the noise level by providing a resistance to the passage of sound emanating from the machine. But, open-cell polyurethane foam has a tendency to soak up oil and the like used to lubricate the machine. Once the cells of the sound-absorbent foam material fill with oil, the material becomes noise reflective, and so is inefficient and ineffective. Further, the accumulated oil represents a fire hazard. Another sound absorbent material is the non-woven fiberglass pad which is similar to open-cell foam in its operation, and likewise becomes ineffective and/or hazardous due to oil absorption.\nOther sound absorbing materials and structures are known in the art, but these materials are too sophisticated, too complicated and/or too expensive. Often the construction materials are expensive, ineffective, inefficient or hazardous, such as with paper, wood or asphalt materials. Such structures are not necessarily inoperative, as they are useful for general acoustical applications, but they are not useful in association with noise generating machinery."}
{"text": "1. Field of the Invention\nThe present invention relates to a hard disk duplication device, especially to a hard disk duplication device capable of expanding duplication capacity.\n2. Description of the Related Art\nConventional hard disk duplication is performed by a PC (personal computer). However, it requires a platform provided by an OS (operating system) software program, and cannot process large number of hard disks at a time.\nBesides, to meet clients' different requirements of different number of hard disks to be processed at a time, the hard disk duplication device manufacturers may have to prepare different specs. For example, assuming the required number of hard disks to be processed at a time for client A is 50 and the required number of hard disks to be processed at a time for client B is 100, then a manufacturer will have to prepare two different specs of hard disk duplication devices. However, to meet the requirements of a variety of specs, the manufacturing cost will increase substantially.\nTo solve the foregoing problem, a novel hard disk duplication device is needed."}
{"text": "In a gas turbine engine fuel delivery system, pump assemblies, as shown for example in US 2005/0232784, are typically used for pumping the fuel. Where such assemblies include gear pump, gear elements are commonly supported by bearing blocks which are adapted to receive respective bearing shafts of the gears through a bore of each bearing block. These bearing blocks also typically abut axially-directed faces of respective gears of the pumps. The bearing blocks may be for solid bearings, or pressure loaded bearings. A solid bearing typically transfers load from journals to the pump housing, and additionally can transfer axial load to the housing. Pressure loaded bearings also transfer load from journals to housing, and in addition can provide an axial force and a moment against the axially-directed face of the gear which the bearing block abuts.\nIt is known to use bimetallic (alloy) bearing blocks, as shown, for example, in U.S. Pat. No. 4,523,365. Such a bimetallic bearing block generally comprises an inner bush covered with an outer backing layer. The inner bush is formed of an alloy which provides a tribologically compatible surface for the gear side face and journal to run against. However it may be undesirable for the whole bearing block to be formed of such an alloy. Two reasons for this are that firstly the weight of a block formed solely of such an alloy may be larger than desired, and secondly the difference between the coefficient of thermal expansion (CTE) and that of the light alloy normally used for the pump housing body may be large. Therefore, in a bimetallic bearing block the inner bush is coated with a light alloy backing layer, which reduces the overall weight of the block, and mitigates the CTE difference with the pump housing body.\nIn order to provide a tribologically compatible surface for the gear side face and journal to run against, an antifriction alloy is typically used to form the inner bush. The antifriction alloy may be, for example, a lead bronze alloy. In particular, the antifriction alloy may be a high concentration lead bronze alloy.\nHowever, testing of bearing blocks which use such an antifriction alloy has shown that such blocks can be prone to suffer permanent radial deformation along the bearing bore when operated above a threshold pressure/temperature combination. If this permanent radial deformation is significant, it can reduce the clearance between the gear journal and the bearing bore. As a result of this reduced clearance, overheating can occur, which may result in mechanical damage of the gear and/or bearing. This problem is more evident in bearing blocks having a higher concentration of lead than those having a lower concentration.\nFIG. 1 shows a typical Stribeck curve, which describes how coefficient of friction varies for different lubrication regimes. It can be seen that typically, higher coefficients of friction occur in a mixed mode lubrication regime as compared to in a full-film lubrication regime. The thrust faces in a gear pump bearing arrangement typically operate at least partly in a mixed-film lubrication regime. This is evident from the wear and scoring that can be visualised at the thrust surfaces after running. In contrast, the journal element of the bearing typically operates in a full-film lubrication regime. Generally, the loads on the thrust face of a bearing block are lower than those acting on the journal element of the bearing, but the sliding velocities are higher. Accordingly, local heating of the thrust face and gear side face is more likely than in the journal element of the bearing. The situation may be particularly acute at the region of the thrust face under the gear root circle diameter where fluid cooling is limited by the restricted fluid flow that occurs across this section of the thrust face. The primary mode for heat transfer away from the thrust face at this point may be conduction through the inner bush and backing layer of the bearing block. Therefore, depending on the thermal conductivity of the alloys of these components, there may be a problem that heat cannot transfer away from the thrust face sufficiently quickly. This may then lead to reduced wear characteristics, and correspondingly poorer performance of the thrust face in a thrust bearing capacity."}
{"text": "A family of toxins, referred to as the Repeats in Toxins (RTX), have been associated with infections caused by gram-negative organisms including E. coli, Proteus, Morganella, Pasteurelia, Actinobacillus, and Bordetella spp. (Menestrina G., et al, Toxicology 87 (1994) 249-267, and references referred to therein). The prototype toxin is Escherichia coli .alpha.-hemolysin which is a primary virulence factor of urinary infections, peritonitis, meningitis, and septicemia caused by virulent strains of E. coli. The RTX toxins share several common features which are discussed in the review by Welch, Mol. Microbiol. 5, 521, 1991.\nThe RTX toxins of Actinobacillus pleuropneumoniae play an important role in infections cause by strains of Actinobacillus pleuropneumoniae (J. Frey et al., in Bacterial Protein Toxins, Zbl.Bakt. Suppl. 24, Freer et al. (Eds)., Gustav Fischer, Stuttgart, Jena, N.Y., 1994). Actinobacillus pleuropneumoniae is the agent responsible for swine pleuropneumonia, a severe contagious disease which causes great economic losses (J. I. MacInnes and N. L Smart, Actinobacillus and Haemophilus in Pathogenesis of Bacterial Infections in Animals, 2nd ed. C. L. Gyles and C. O. Thoen, Iowa State University Press, Ames, Chapter 16). A. pleuropneumoniae is a gram-negative bacterium of the family Pasteurellaceae J. I. MacInnes and N. L Smart, supra). Twelve serotypes of the organism have been described by serotyping based on capsular polysaccharides (Nicolet, J. Can. Vet. J. 29:578-580, 1988 and Nielsen, R., cta. Vet. Scand. 27:453-455, 1986). Differences in virulence have been observed among the serotypes and a number of virulence factors have been considered responsible for the differences. These factors include capsular polysaccharides, liposaccharides, exotoxins and adhesion factors (J. I. MacInnes and N. L Smart, supra).\nTwo hemolytic RTX toxins, ApxI and ApXII, and one non-hemolytic RTX toxin, ApxIII have been identified in different A. pleuropneumoniae serotype strains J. Frey et al., supra) The toxins are secreted into the growth medium by A. pleuropneumoniae, they have different molecular masses, and they can be distinguished serologically by polyclonal and monoclonal antibodies (J. Frey et al., supra). It has been reported that A. pleuropneumoniae serotypes which produce a combination of two Apx toxins are more virulent than serotypes with one toxin alone (J. Frey et al., supra). The serotypes producing ApxI and ApxII are the most virulent (J. Frey et al., supra). No A. pleuropneumoniae strains have been identified which produce all three toxins J. Frey et al., supra).\nRTX toxins have also been reported to be secreted in other bacterial strains including Bordetella pertussis (Betsou et al.1993 and CA 1,189,790) E. coli, Listeria, Moraxella, Pseudomonas, Staphylococcus, Vibrio (Nakai et al. 1983) and Neisseria meningitidis (Thompson et al.1993).\nVaccines which have been developed for preventing infections by A. pleuropneumoniae have been based on whole live cells, attenuated cells, lysates, culture supernatants, and extracts of A. pleuropneumoniae. Canadian Patent 1,189,790 describes a vaccine containing A. pleuropneumoniae cells, cell fragments, extracts and/or metabolites, and an adjuvant derived from Bordetella pertussis. Other proposed vaccines contain: (a) inactivated toxin of serotype I and optionally a toxin of another serotype (EP-A-420.743); (b) whole cell or sonicated whole cell components of a virulent strain obtained after passage in a host (WO9321951-A); (c) transferring binding protein, cytolysin and/or APP4 (WO9308283-A); (d) at least one immunogenic part of at least one cytolytic A. pleuropneumoniae protein prepared by recombinant DNA methods (Canadian Patent Application 2,045,950); (e) an iron-repressible outer membrane protein of molecular weight 105 kD (Canadian Patent Application 2,045,950); (f) outer membrane proteins having a major dominant antigenic protein component of 42 kD and a haemolysin of 105 kD and/or macrophage toxin of 120 kD (Canadian Patent Application No. 2,040,544); (g) inactivated toxin of serotype 1 of A. pleuropneumoniae which is obtained from culture supernatant (EP-420743); (h) extracellular proteinaceous materials from the culture medium of strains of at least two different serotypes of A. pleuropneumoniae. (EP-420743); or, (i) hemolysin antigen produced by recombinant techniques. Many of the known vaccines have limited effectiveness particularly against infection by heterologous serotypes.\nConventional formalin-killed bacterins provide limited protection against challenge with homologous serotypes of A. pleuropneumoniae and poor protection against heterologous serotypes (Neilson, 1984). In contrast, convalescent pigs are completely protected from challenge with homologous serotypes and significantly protected from disease from heterologous serotypes (Neilson, 1984).\nMany of the antigenic components of A. pleuropneumoniae which have been investigated as potential vaccine candidates all fall short of affording complete cross-protection. Devenish et al. (1990) demonstrated homologous protection from challenge when gel-purified ApxI and ApxII cytolysins were used as a vaccine. However, others have reported less than complete protection with partially purified toxin vaccines, although the protection afforded is increased when these vaccines are enriched with other cellular components such as outer membrane proteins (OMPs) (Van den Bosch et al., 1992).\nLive attenuated vaccines made from strains of A. pleuropneumoniae lacking in capsule production have been reported to protect against homologous challenge (Rosendal et al., 1990; Inzana et al., 1988). Results have not been reported with heterologous challenge. Live strains deficient in RTX toxin production have been reported to afford no protection against disease (Inzana et al., 1991). Jansen (1994) has concluded that both opsonization of A. pleuropneumoniae to enhance phagocytosis and neutralization of the RTX toxins are necessary for immune protection against disease."}
{"text": "With more and more functions integrated within terminal devices such as mobile phones, it is more convenient for users to use terminal devices due to the integrated new functions.\nAll major manufacturers of terminal devices bring their attention to fingerprint identification function recently. For conventional terminal devices such as mobile phones, tablet computers, when designing the fingerprint identification function, the fingerprint identification function is generally combined with a physical button, a section view of which is shown in FIG. 1, so as to make the fingerprint identification function compatible with the button. Generally, with such design, it is required to form an opening in a fingerprint identification area of a glass cover 01, because a physical button 02 can only be activated through a fixed travel of press. In this way, when a fingerprint identification sensor 03 is tapped by a user, only a fingerprint may be identified by the system. Moreover, when the fingerprint identification sensor 03 is pressed through a certain travel by the user, the physical button 02 may be activated, and two actions including fingerprint sensing and button-press may be identified by the system. Then, a corresponding action may be extracted according to decisions of a upper level processor in the terminal device. However, the exterior appearance of the terminal device is affected by the two-layers structure described above, and the user operation become inconvenient as well."}
{"text": "The present invention relates to the field of the transmission of coded audio signals, more specifically to a method and an apparatus for obtaining, or acquiring, spectrum coefficients for a replacement frame of an audio signal, to an audio decoder, to an audio receiver and to a system for transmitting audio signals. Embodiments relate to an approach for constructing a spectrum for a replacement frame based on previously received frames.\nIn conventional technology, several approaches are described dealing with a frame-loss at an audio receiver. For example, when a frame is lost on the receiver side of an audio or speech codec, simple methods for the frame-loss-concealment as described in P. Lauber and R. Sperschneider, “Error Concealment for Compressed Digital Audio,” in AES 111th Convention, New York, USA, 2001 (hereinafter “the Lauber reference”) may be used, such as:\nrepeating the last received frame,\nmuting the lost frame, or\nsign scrambling.\nAdditionally, in the Lauber reference, an advanced technique using predictors in sub-bands is presented. The predictor technique is then combined with sign scrambling, and the prediction gain is used as a sub-band wise decision criterion to determine which method will be used for the spectral coefficients of this sub-band.\nIn U.S. Pat. No. 6,351,730 B2 (C. J. Hwey, “Low-complexity, low-delay, scalable and embedded speech and audio coding with adaptive frame loss concealment,” hereinafter “the '730 Patent”), a waveform signal extrapolation in the time domain is used for a MDCT (Modified Discrete Cosine Transform) domain codec. This kind of approach may be good for monophonic signals including speech.\nIf one frame delay is allowed, an interpolation of the surrounding frames can be used for the construction of the lost frame. Such an approach is described in US Patent Application Publication No. 2007/094009 A1 (S. K. Gupta, E. Choy and S.-U. Ryu, “Encoder-assisted frame loss concealment techniques for audio coding,” hereinafter “the '009 Publication”), where the magnitudes of the tonal components in the lost frame with an index m are interpolated using the neighboring frames indexed m−1 and m+1. The side information that defines the MDCT coefficient signs for tonal components is transmitted in the bit-stream. Sign scrambling is used for other non-tonal MDCT coefficients. The tonal components are determined as a predetermined fixed number of spectral coefficients with the highest magnitudes. This approach selects n spectral coefficients with the highest magnitudes as the tonal components.Cm*(k)=½(Cm−1(k)+) Cm+1(k)\nFIG. 7 shows a block diagram representing an interpolation approach without transmitted side information as it is for example described in S.-U. Ryu and K. Rose, “A Frame Loss Concealment Technique for MPEG-AAC,” in 120th AES Convention, Paris, France, 2006 (hereinafter “Ryu 2006/Paris”. The interpolation approach operates on the basis of audio frames coded in the frequency domain using MDCT (modified discrete cosine transform). A frame interpolation block 700 receives the MDCT coefficients of a frame preceding the lost frame and a frame following the lost frame, more specifically in the approach described with regard to FIG. 7, the MDCT coefficients Cm−1(k) of the preceding frame and the MDCT coefficients Cm+1(k) of the following frame are received at the frame interpolation block 700. The frame interpolation block 700 generates an interpolated MDCT coefficient Cm(k) for the current frame which has either been lost at the receiver or cannot be processed at the receiver for other reasons, for example due to errors in the received data or the like. The interpolated MDCT coefficient Cm (k) output by the frame interpolation block 700 is applied to block 702 causing a magnitude scaling in scale factor band and to block 704 causing a magnitude scaling with an index set, and the respective blocks 702 and 704 output the MDCT coefficient Cm(k) scaled by the factor â(k) and ã(k), respectively. The output signal of block 702 is input into the pseudo spectrum block 706 generating on the basis of the received input signal the pseudo spectrum {circumflex over (P)}m(k) that is input into the peak detection block 708 a signal indicating detected peaks. The signal provided by block 702 is also applied to the random sign change block 712 which, responsive to the peak detection signal generated by block 708, causes a sign change of the received signal and outputs a modified MDCT coefficient Ĉm (k) to the spectrum composition block 710. The scaled signal provided by block 704 is applied to a sign correction block 714 causing, in response to the peak detection signal provided by block 708 a sign correction of the scaled signal provided by block 704 and outputting a modified MDCT coefficient {tilde over (C)}m(k) to the spectrum composition block 710 which, on the basis of the received signals, generates the interpolated MDCT coefficient Cm*(k) that is output by the spectrum composition block 710. As is shown in FIG. 7, the peak detection signal provided by block 708 is also provided to block 704 generating the scaled MDCT coefficient.\nFIG. 7 generates at the output of the block 714 the spectral coefficients {tilde over (C)}m(k) for the lost frame associated with tonal components, and at the output of the block 712 the spectral coefficients Ĉm(k) for non-tonal components are provided so that at the spectrum composition block 710 on the basis of the spectral coefficients received for the tonal and non-tonal components the spectral coefficients for the spectrum associated with the lost frame are provided.\nThe operation of the FLC (Frame Loss Concealment) technique described in the block diagram of FIG. 7 will now be described in further detail.\nIn FIG. 7, basically, four modules can be distinguished:                a shaped-noise insertion module (including the frame interpolation 700, the magnitude scaling within the scale factor band 702 and the random sign change 712)        a MDCT bin classification module (including the pseudo spectrum 706 and the peak detection 708),        a tonal concealment operations module (including the magnitude scaling within the index set 704 and the sign correction 714), and        the spectrum composition 710.        \nThe approach is based on the following general formula:Cm(k)=Cm*(k)a*(k)s*(k), 0≤k<M \nCm*(k) is derived by a bin-wise interpolation (see block 700 “Frame Interpolation”):Cm*(k)=½(Cm−1(k)+Cm+1(k))\nα*(k) is derived by an energy interpolation using the geometric mean:                scale factor band wise for all components, (see block 702 “Magnitude Scaling in Scalefactor Band”) and        index sub-set wise for tonal components (see block 704 “Magnitude Scaling within Index Set”):        \n                    (                  α          *                )            2        ⁢          (      k      )        =                              E                      m            +            1                          ⁢                  E                      m            -            1                                      E      m      \nFor tonal components it can be shown that a=cos (πfl), with fl being the frequency of the tonal component.\nThe energies E are derived based on a pseudo power spectrum, derived by a simple smoothing operation:P(k)≅C2(k)+{C(k+1)−C(k−1)}2 \ns*(k) is set randomly to ±1 for non-tonal components (see block 712 “Random Sign Change”), and to either +1 or −1 for tonal components (see block 714 “Sign Correction”).\nThe peak detection is performed as searching for local maxima in the pseudo power spectrum to detect the exact positions of the spectral peaks corresponding to the underlying sinusoids. It is based on the tone identification process adopted in the MPEG-1 psychoacoustic model described in ISO/IEC JTC1/SC29/WG11, Information technology—Coding of moving pictures and associated, International Organization for Standardization, 1993. Out of this, an index sub-set is defined having the bandwidth of an analysis window's main-lobe in terms of MDCT bins and the detected peak in its center. Those bins are treated as tone dominant MDCT bins of a sinusoid, and the index sub-set is treated as an individual tonal component.\nThe sign correction s*(k) flips either the signs of all bins of a certain tonal component, or none. The determination is performed using an analysis by synthesis, i.e., the SFM is derived for both versions and the version with the lower SFM is chosen. For the SFM derivation, the power spectrum is needed, which in return may use the MDST (Modified Discrete Sine Transform) coefficients. For keeping the complexity manageable, only the MDST coefficients for the tonal component are derived, using also only the MDCT coefficients of this tonal component.\nFIG. 8 shows a block diagram of an overall FLC technique which, when compared to the approach of FIG. 7, is refined and which is described in S.-U. Ryu and R. Kenneth, An MDCT domain frame-loss concealment technique for MPEG Advanced Audio Coding, Department od Electrical and Computer Engineering, University of California, 2007 (hereinafter “Ryu 2007”). In FIG. 8, the MDCT coefficients Cm−1 and Cm+1 of a last frame preceding the lost frame and a first frame following the lost frame are received at an MDCT bin classification block 800. These coefficients are also provided to the shape-noise insertion block 802 and to the MDCT estimation for a tonal components block 804. At block 804 also the output signal provided by the classification block 800 is received as well as the MDCT coefficients Cm−2 and Cm+2 of the second to last frame preceding the lost frame and the second frame following the lost frame, respectively, are received. The block 804 generates the MDCT coefficients {tilde over (C)}m of the lost frame for the tonal components, and the shape-noise insertion block 802 generates the MDCT spectral coefficients for the lost frame Ĉm for non-tonal components. These coefficients are supplied to the spectrum composition block 806 generating at the output the spectral coefficients Cm* for the lost frame. The shape-noise insertion block 802 operates in reply to the system IT generated by the estimation block 804.\nThe following modifications are of interest with respect to the Ryu 2006/Paris reference:                The pseudo power spectrum used for the peak detection is derived asPm(k)=Cm−12(k)+Cm+12(k)        To eliminate perceptually irrelevant or spurious peaks, the peak detection is only applied to a limited spectral range and only local maxima that exceed a relative threshold to the absolute maximum of the pseudo power spectrum are considered. The remaining peaks are sorted in descending order of their magnitude, and a pre-specified number of top-ranking maxima are classified as tonal peaks.        The approach is based on the following general formula (with a being signed this time):Cm(k)=Cm*(k)α(k), 0≤k<M         Cm*(k) is derived as above, but the derivation of a becomes more advanced, following the approachEm(α)=½{Em−1(α)+Em+1(α)}        Substituting Em, Em−1, and Em+1 withEm−1(α)≅|cm−1|2+|sm−1|2=|cm−1|2+|ξ1+αζ1|2 Em(α)≅α2|cm|2+|sm|2+α2|cm|2+|ξ2+αζ2|2 Em+1(α)≅|cm+1|2+|sm+1|2=|cm+1|2+|ξ3+αζ3|2 whereassm−1≅A1cm−2+A2cm−1+αA3cm=ξ1+αζ1 sm≅A1cm−1αlA2cm+A3cm+1=ξ2+αζ2 sm+1≅αA1cm+A2cm+1+A3cm+2=ξ3+αζ3         yields an expression that is quadratic in α. Hence, for the given MDCT estimate there exist two candidates (with opposite signs) for the multiplicative correction factor (A1, A2, A3 are the transformation matrices). The selection of the better estimate is performed similar to what is described in the Ryu 2006/Paris reference.        This advanced approach may use two frames before and after the frame loss in order to derive the MDST coefficients of the previous and the subsequent frame.        \nA delay-less version of this approach is suggested in S.-U. Ryu, Source Modeling Approaches to Enhanced Decoding in Lossy Audio Compression and Communication, UNIVERSITY of CALIFORNIA Santa Barbara, 2006 (hereinafter “Ryu 2006/California”):                As a starting point, the interpolation formula Cm*(k)=½(Cm−1(k)+Cm+1(k)) is reused, but is applied for the frame m−1, resulting in:Cm(k)=2Cm−1*(k)−Cm−2(k)        \nThen, the interpolation result Cm−1* is replaced by the true estimation (here, the factor 2 becomes part of the correction factor: α=2 cos (πfl)), which leads toCm(k)=αCm−1(k)−Cm−2(k)\nThe correction factor is determined by observing the energies of two previous frames. From the energy computation, the MDST coefficients of the previous frame are approximated assm−1≅(A1−A3)cm−2+A2cm−1+αA3cm−1=ξ0+αζ0                 Then, the sinusoidal energy is computed asEm−1(α)≅|cm−1|2+|sm−1|2=|cm−1|2+|ξ0+αζ0|2         Similarly, the sinusoidal energy for frame m−2 is computed and denoted by Em−2, which is independent of a.        Employing the energy requirementEm−1(α)=Em−2         yields again an expression that is quadratic in a.        The selection process for the candidates computed is performed as before, but the decision rule accounts only the power spectrum of the previous frame.        \nAnother delay-less frame-loss-concealment in the frequency domain is described in European Patent No. EP 0574288 B1 (M. Yannick, “Method and apparatus for transmission error concealment of frequency transform coded digital audio signals,” hereinafter “the '288 Patent”. The teachings of reference the '288 Patent can be simplified, without loss of generality, as:                Prediction using a DFT of a time signal:                    (a) Obtain the DFT spectrum from the decoded time domain signal that corresponds to the received coded frequency domain coefficients Cm.            (b) Modulate the DFT magnitudes, assuming a linear phase change, to predict the missing frequency domain coefficients in the next frame Cm+1.                        Prediction using a magnitude estimation from the received frequency spectra:                    (a) Find C′m and S′m, using Cm as input, such thatC′m(k)=Qm(k)cos(φm(k)+χ)S′m(k)=Qm(k)sin(φm(k)+χ)            where Qm(k) is the magnitude of the DFT coefficient that corresponds to Cm(k).            (b) Calculate:                        \n                    Q        m            ⁡              (        k        )              =                                                                                  C                m                ′                            ⁡                              (                k                )                                                          2                +                                                                        S                m                ′                            ⁡                              (                k                )                                                          2                                        φ        m            ⁡              (        k        )              =          arccos      ⁢                                    C            m                    ⁡                      (            k            )                                                Q            m                    ⁡                      (            k            )                                                                  (c) Perform a linear extrapolation of the magnitude and the phase:Qm+1(k)=2Qm(k)−Qm−1(k)φm+1(k)=2φm(k)−φm−1(k)Cm+1(k)=Qm+1(k)cos(φm+1(k))                        Use filters to calculate C′m and S′m from Cm and then proceed as above to get Cm+1(k)        Use an adaptive filter to calculate Cm+1(k):        \n            C              m        +        1              ⁡          (      k      )        =                    ∑                  i          =          0                I            ⁢                        a                      m            ,            i                          ⁡                  (          k          )                      +                  C                  m          -          i                    ⁡              (        k        )            \nThe selection of spectrum coefficients to be predicted is mentioned in the '288 Patent but is not described in detail.\nIn Y. Mahieux, J.-P. Petit and A. Charbonnier, “Transform coding of audio signals using correlation between successive transform blocks,” in Acoustics, Speech, and Signal Processing, 1989. ICASSP-89., 1989, it has been recognized that, for quasi-stationary signals, the phase difference between successive frames is almost constant and depends only on the fractional frequency. However, only a linear extrapolation from the last two complex spectra is used.\nIn AMR-WB+ (see 3GPP; Technical Specification Group Services and System Aspects, Extended Adaptive Multi-Rate—Wideband (AMR-WB+) codec, 2009) a method described in U.S. Pat. No. 7,356,748 B2 (A. Taleb, “Partial Spectral Loss Concealment in Transform Codecs,” hereinafter “the '748 Patent”) is used. The method in the '748 Patent is an extension of the method described in reference the '288 Patent in a sense that it uses also the available spectral coefficients of the current frame, assuming that only a part of the current frame is lost. However, the situation of a complete loss of a frame is not considered in the '748 Patent.\nAnother delay-less frame-loss-concealment in the MDCT domain is described in US Patent Application Publication No. 2012/109659 A1 (C. Guoming, D. Zheng, H. Yuan, J. Li, J. Lu, K. Liu, K. Peng, L. Zhibin, M. Wu and Q. Xiaojun, “Compensator and Compensation Method for Audio Frame Loss in Modified Discrete Cosine Transform Domain,” hereinafter “the '659 Publication”. In the '659 Publication, it is first determined if the lost Pth frame is a multiple-harmonic frame. The lost Pth frame is a multiple-harmonic frame if more than K0 frames among K frames before the Pth frame have a spectrum flatness smaller than a threshold value. If the lost Pth frame is a multiple-harmonic frame then (P−K)th to (P−2)nd frames in the MDCT-MDST domain are used to predict the lost Pth frame. A spectral coefficient is a peak if its power spectrum is bigger than the two adjacent power spectrum coefficients. A pseudo spectrum as described in L. S. M. Dauder, “MDCT Analysis of Sinusoids: Exact Results and Applications to Coding Artifacts Reduction,” IEEE TRANSACTIONS ON SPEECH AND AUDIO PROCESSING, pp. 302-312, 2004 (hereinafter “Dauder”), is used for the (P−1)st frame.\nA set of spectral coefficients Sc is constructed from L1 power spectrum frames as follows.\nObtaining L1 sets S1, . . . , SL1 composed of peaks in each of L1 frames, a number of peaks in each set being N1, . . . , NL1, respectively. Selecting a set Si from the L1 sets of S1, . . . , SL1. For each peak coefficient mj, j=1 . . . Ni in the set S1, judging whether there is any frequency coefficient among mj, mj±1, . . . , mj±k belonging to all other peak sets. If there is any, putting all the frequencies mj, mj±1, . . . , mj±k into the frequency set SC. If there is no frequency coefficient belonging to all other peak sets, directly putting all the frequency coefficients in a frame into the frequency set SC. Said k is a nonnegative integer. For all spectral coefficients in the set SC the phase is predicted using L2 frames among (P−K)th to (P−2)nd MDCT-MDST frames. The prediction is done using a linear extrapolation (when L2=2) or a linear fit (when L2>2). For the linear extrapolation:\n                    φ        ^            p        ⁡          (      m      )        =                    φ                  t          ⁢                                          ⁢          1                    ⁡              (        m        )              +                            p          -                      t            ⁢                                                  ⁢            1                                                t            ⁢                                                  ⁢            1                    -                      t            ⁢                                                  ⁢            2                              ⁡              [                                            φ                              t                ⁢                                                                  ⁢                1                                      ⁡                          (              m              )                                -                                    φ                              t                ⁢                                                                  ⁢                2                                      ⁡                          (              m              )                                      ]            where p, t1 and t2 are frame indices.\nThe spectral coefficients not in the set SC are obtained using a plurality of frames before the (P−1)st frame, without specifically explaining how."}
{"text": "This invention relates to a new and improved method for refining coal, particularly coal containing sulfur.\nVarious methods have been proposed for removal of sulfur from coal. The methods are generally tailored to remove a particular form of sulfur which may be an elemental, pyritic, organic or sulfate form of sulfur. Because of the extremely complex composition of coal and the radical variance in composition from seam to seam, and even within a single seam, the problems in effectively removing substantially all sulfur and sulfur compounds are significant.\nIt is believed, that the goal of reducing sulfur emissions from the combustion of coal can best be achieved prior to pyrolysis using efficient refining methods. In this manner combustion systems can be designed to a maximized efficiency through fueling with a uniform, high Btu, low ash and low-sulfur feed stock. Use of efficient combustion systems and high grade refined coal will have a substantial impact on acid raid caused by sulfur dioxide emissions from coal-fired energy and processing plants.\nIt is a primary object of this invention to produce a moisture free, fine particle material that can be used, as is, in a blow-fired or fluidized bed combustor; formed into a briquette or slug, for a grate-type combustor; or, mixed in an slurry or emulsion for transport and firing in a liquid spray as an oil substitute.\nIt has become apparent from technology trends that future transportation methods and combustion methods will utilize finely ground coal, often below 200 mesh. In certain instances where emulsions with oil or water or mixtures thereof are formulated, coal is ground to powders 350-400 mesh and below. While it may be advantageous to finely grind coal to this size to liberate the carbonaceous materials from the inert ash and chemically undesirable compounds as pyrite, such finely ground particles present substantial problems in applying conventional separation procedures.\nIn conventional systems for cleaning coal, it is customary to remove substantial portions of the pyrites and ash using largely mechanical separation means. Coal is crushed to a uniform size and specific gravity techniques are used to separate the lighter weight carbon and hydrocarbon materials from the heavier pyrite and ash waste materials. One effective method of separation is with a water cyclone. A water cyclone or hydrocyclone is a concentrator for the wet processing, of finely divided materials. The cyclone has no moving parts, and utilizes the principles of centrifugal forces to separate materials of different specific gravity by directing a high velocity slurry at a steep surface of a conical bowl, which provides for an overflow carrying lighter weight materials, and an underflow carrying heavier weight materials. By careful design of the geometry of the surface, materials can be separated at select median specific gravities or cut points. For coal, complex multiple angle surfaces have been devised to provide a full range of selectable cut points adaptable to the variety of different coals. Such compound water cyclones are customerily used in multiples with a recirculating of middling material to improve a clear separation of lights and heavies in desired proportions.\nThe composition of coal renders the specific gravity separation of constituent materials an imprecise art. Coal is an amalgamation of different materials which is not consistently amorphous and is not strictly stratified. Liberating particles of the constituent materials by crushing improves the potential ability to separate at desired specific gravities. However, it has been found that there is a limit for grinding using conventional cyclone systems because fine particles in the range of 200 mesh and below tend to naturally agglomerate in the cyclone feedwater system, thereby inhibiting particle differentiation and select separation in the cyclone. While the cause of such balling is not wholly understood, it is believed that the substantial increase in surface area in a mass of extremely fine particles allow exposed hydrophobic hydrocarbons of the carbonaceous material, which have a natural mutual affinity, to cause a coalescence of particles into an agglomerate. While the use of selective surfactants can encourage this agglomeration process and lead to methods for select separation of carbonaceous and non-carbonaceous compositions, such methods are somewhat slower than cycloning yet are recognized as competitive processes.\nThe method of refining coal here devised solves the problem of undersirable agglomeration and provides additional benefits improving the resultant efficiency of the refining operation and the quality of the refined product."}
{"text": "Electrochemical sensor elements are used in laboratory and process measurement technology in many areas of chemistry, biochemistry, pharmacy, biotechnology, food technology, water management, and environmental monitoring for the analysis of measuring media—in particular, of measuring liquids. Electrochemical measuring techniques allow detection of, for example, activities of chemical substances, such as ions, and of a therewith-correlated process variable of an—in particular, liquid—measuring medium. The substance whose concentration or activity is to be measured is also referred to as analyte. Electrochemical sensor elements may, for example, be potentiometric or amperometric sensor elements.\nPotentiometric sensor elements typically comprise a measuring half-cell and a reference half-cell, as well as a measurement circuit. The measuring half-cell includes a potential-forming element, which may comprise, for example, a redox electrode, an analyte-sensitive coating, or an ion-selective membrane, depending upon the type of potentiometric sensor element. The measurement circuit generates a measurement signal that represents the potential difference between the measuring half-cell and the reference half-cell. Examples of potentiometric sensor elements are ion-selective electrodes (ISE). A special case of such an ion-selective electrode is the known pH glass electrode, which has a chamber which is closed by a pH-sensitive glass membrane and includes an internal electrolyte contained therein and a discharge element immersed therein.\nAmperometric sensor elements, e.g., amperometric disinfection sensors, may, for example, have an—in particular, potentiostatic—control circuit which is designed to preset a nominal voltage between a working electrode and a reference electrode and to detect the current flowing between the working electrode and the counter electrode in the process. One example of an amperometric sensor element is described in DE 10 2008 039 465 A1, for example.\nOther examples of electrochemical sensor elements are such as are based upon electrolyte isolator semiconductor layer stacks (for short: EIS). These layer systems may, for example, be operated capacitively or as ion-selective field-effect transistors (ISFET). DE 198 57 953 A1 describes such a sensor element designed to measure the ion concentration or the pH value of a measuring medium using an ISFET.\nWithin the scope of this application, the term, “electrochemical sensor elements,” furthermore includes conductively- or capacitively-operating conductivity sensors, as well as (spectro)photometrically-operating sensors, such as turbidity sensors.\nThe sensor elements are in this case in contact with the process via the process connection and thus typically, at least partially and at least intermittently or, in particular, substantially mostly, come into contact with the measuring medium. In the case where the sensor element is an electrochemical sensor element used in the aforementioned industries, regular removal of the sensor element from the process is necessary.\nOn the one hand, the removal serves to clean the sensor element at regular intervals—for example, in order to cleanse it of bacteria and other microorganisms. In this respect, autoclaving, in which the sensor element is cleaned at high temperature (approx. 140° C.) and high pressure using water vapor, has proven itself as the sterilization method. Autoclaving is used, in particular, for electrochemical sensor elements—for example, in the pharmaceutical industry, food industry, medical technology, and to her fields.\nOn the other hand, calibration, verification, and/or adjustment of the sensor element may be necessary at regular intervals. In this case, calibration usually means the determination of a deviation of the first measured value measured by the sensor element from a reference measured value provided by a second sensor element and accepted as correct. The verification also comprises the determination of the deviation and its assessment or evaluation. Lastly, adjustment means the adaptation of the sensor element in such a way that a model used by the sensor element to determine a measured value from a raw value is adapted such that the measured value corresponds to the reference value provided by the second sensor element. The sensor element also has to be removed from the process for calibration, verification, and/or adjustment.\nDepending upon the degree to which the electrochemical sensor element is stressed and/or damaged in the process by the contact with the measuring medium and/or the aforementioned regular cleaning, a regular exchange of the sensor element may also be necessary.\nIn order to process and/or transmit electronic and/or electrical signals generated by the sensor element—for example, to a transmitter unit or a higher-level unit, e.g., a control center—the sensor element is often connected to a connection element.\nThe connection element comprises a cable attachment and a cable, which in turn is connected to the higher-level unit. An interface that is, for example, designed to be inductive or optical is respectively located on the sensor element and on the connection element. Via these interfaces, the sensor element is supplied with energy, and/or communication between the sensor element and the connection element is ensured; this is described, for example, in EP 1 625 643.\nIn this respect, reference is also made, in particular, to the products with the name, “Memosens,” by the applicant. Other generic designs are, for example, “Memosens” by the company Knick, “ISM” by Mettler-Toledo, the “ARC” system by Hamilton, and “SMARTSENS” by Krohne.\nThe great advantage of the known “Memosens” plug connection is that the above-described regular removal of the sensor element from the process is easily possible. The sensor element can be removed from the process easily and can, where applicable, be brought into the laboratory for calibration and/or adjustment (for example, by means of a so-called “Memobase Plus” system). When a first sensor element is removed, for example, it can, moreover, be replaced directly at the measuring site close to the process by a second sensor element. As a result, the maintenance duration or idle times of the process system and, ultimately, also the costs are drastically reduced. Another advantage is that the electronic components of the sensor element and of the connection element are protected as a result of their—where applicable—liquid-tight design when the plug connection is disconnected. This is important in order to prevent, for example, moisture ingress during a hot steam jet cleaning used for hygienic reasons in the industries mentioned at the beginning.\nA disadvantage of the connection by means of the connection element with cable attachment and cable is that a connected transmitter unit must be placed at a distance from the sensor element itself. An arrangement of the transmitter unit close to the process or close to the sensor element is, however, often desired, even if only to be able to calculate and display the process variable, which can be determined by the sensor element in combination with the transmitter unit, close to the process. Within the scope of the application, a “transmitter unit” refers to a unit which comprises electronic components and by means of which at least the process variable of the sensor element can be determined completely, and with which a standardized transmission of the process variable to a higher-level unit is possible. For this purpose, the transmitter unit has, for example, a microprocessor and/or a storage unit.\nSuch a standardized transmission takes place using a standardized wired or wireless communications network. In the simplest case, this network is formed as a measurement transmission path according to the 4-20 mA standard. A wired communications network can, for example, furthermore be a wired fieldbus of automation technology, such as Foundation Fieldbus, Profibus PA, Profibus DP, HART, CANBus, etc. It can, however, also be a modern industrial communications network, such as an “Industrial Ethernet” fieldbus—in particular, Profinet, HART-IP, or Ethernet/IP, or a communications network known from the communications sector, such as Ethernet according to the TCP/IP protocol. In the case of a wireless communications network, it can, for example, be a Bluetooth, ZigBee, WLAN, GSM, LTE, UMTS communications network, or even a wireless version of a fieldbus—in particular, of a standard based upon 802.15.4, such as WirelessHART."}
{"text": "The present invention relates to an array antenna transmitter/receiver which suppresses interference with another user by control of antenna directivity and, more particularly, to a multi-beam antenna transmitter/receiver, transmitting/receiving method, and transmission beam selection method which select transmission/reception directivity from a plurality of fixed directional patterns (multi beams).\nIn a cellular mobile communication system and the like, it has been examined to apply a method of forming a directional pattern (beam) which increases the transmission/reception gain in a desired signal direction and decreases it in another direction by using an array antenna transmitter/receiver including a plurality of antenna elements for the purpose of a higher-speed, higher-quality signal and a larger subscriber capacity. One of such methods is a multi-beam method of selecting a transmission/reception beam from a plurality of fixed directional patterns (multi beams).\nAs disclosed in, e.g., “Multi-Beam Antenna System for Radio Base Station” (Japanese Patent Laid-Open No. 11-266228), a multi-beam antenna transmitter/receiver of this type selects and receives a reception beam having a delay path excellent in reception quality from a plurality of fixed reception beams in reception. In transmission, the multi-beam antenna transmitter/receiver selects and transmits a transmission beam in the same direction as that of a path delay/reception beam number pair excellent in reception quality from path delay/reception beam number pairs selected upon reception.\nFIG. 7 is a block diagram showing an example of a conventional multi-beam antenna transmitter/receiver. The conventional multi-beam antenna transmitter/receiver comprises a reception array antenna 201, an antenna 1 radio reception unit 2031 to antenna N radio reception unit 203N corresponding to reception antenna elements 2021 to 202N, a reception beam 1 formation unit 2041 to reception beam M formation unit 204M (to be also referred to as reception beam formation units 204), a user 1 demodulation block 2051 to user L demodulation block 205L (to be also referred to as user demodulation blocks 205), a user 1 modulation unit 2111 to user L modulation unit 211L, a user 1 transmission beam switching circuit 2121 to user L transmission beam switching circuit 212L, a transmission beam 1 formation unit 2131 to transmission beam J formation unit 213J, an antenna 1 radio transmission unit 2141 to antenna K radio transmission unit 214K corresponding to transmission antenna elements 2161 to 216K, and a transmission array antenna 215.\nThe reception array antenna 201 is formed from the N reception antenna elements 2021 to 202N. Each of the reception antenna elements 2021 to 202N is free from any limitation on the horizontal and vertical directivities, and has, for example, omnidirectivity or dipole. The N reception antenna elements 2021 to 202N are arranged close to each other so as to correlate the reception signals of the antenna elements with each other. The reception array antenna 201 is not limited in the number of reception antenna elements and their layout as far as the N reception antenna elements 2021 to 202N are arranged close to each other. An example of the layout is a circular layout or linear layout at the half-wavelength interval of the carrier.\nSignals received by the N reception antenna elements 2021 to 202N contain desired user signal components, interference signal components, and thermal noise. Each of the desired user signal component and interference signal component contains multipath components. In general, these signal components (desired user signal component and interference signal component) arrive from different directions. Thus, pairs of path delays and reception beam numbers (path delays/reception beam numbers) of desired user signals exist.\nEach of the antenna 1 radio reception unit 2031 to antenna N radio reception unit 203N comprises a low-noise amplifier, bandpass filter, mixer, local oscillator, AGC (Auto Gain Controller), quadrature detector, low-pass filter, analog/digital converter, and the like. The antenna 1 radio reception unit 2031 will be taken as an example. The antenna 1 radio reception unit 2031 receives an output from the reception antenna element 2021, performs reception processes such as amplification of an input signal, frequency conversion from the radio band to the baseband, quadrature detection, and analog/digital conversion, and outputs the resultant signal to the reception beam 1 formation unit 2041 to reception beam M formation unit 204M.\nThe reception beam 1 formation unit 2041 to reception beam M formation unit 204M receive outputs from the antenna 1 radio reception unit 2031 to antenna N radio reception unit 203N, form fixed reception beams different between the respective reception beam formation units for the input signals, and output the beams to the user 1 demodulation block 2051 to user L demodulation block 205L. The number of fixed reception beams, the shape of the fixed reception beam, and the fixed reception beam formation method are not specifically limited. An example of the shape of the fixed reception beam is a quadrature multi-beam, and an example of the fixed reception beam formation method is a method of multiplying input signals by a fixed complex beam weight by digital calculation and calculating the sum.(digital beam forming). In FIG. 7, the reception beam 1 formation unit 2041 to reception beam M formation unit 204M are arranged on the output side of the antenna 1 radio reception unit 2031 to antenna N radio reception unit 203N, and form beams for digital signals of the baseband. A beam formation method in the radio band such as Butler matrix can also be adopted.\nThe reception beam 1 formation unit 2041 to reception beam M formation unit 204M form fixed reception beams different between the respective reception beam formation units 204 for input signals containing the components of all user signals (user 1 signal to user L signal) and the multipath components of the user signals, and demultiplex the input signals in respective arrival directions.\nEach of the user 1 demodulation block 2051 to user L demodulation block 205L is formed from a reception beam 1 path detection unit 2061 to reception beam M path detection unit 206M, a path delay/reception beam selection unit 207, a transmission beam selection unit 209, and a demodulation unit 210.\nThe user 1 demodulation block 2051 to user L demodulation block 205L output user 1 reception data to user L reception data (user reception data) in correspondence with respective users. Since the user demodulation blocks 205 have the same function, the user 1 demodulation block 2051 will be exemplified.\nThe user 1 demodulation block 2051 receives outputs from the reception beam 1 formation unit 2041 to reception beam M formation unit 204M, and outputs a user 1 transmission beam number and user 1 reception data.\nThe reception beam 1 path detection unit 2061 to reception beam M path detection unit 206M receive outputs from the reception beam 1 formation unit 2041 to reception beam M formation unit 204M, detect the path delays of user signals in the input signals, measure the reception qualities of the user signals at the detected path delays, and output reception quality information to the path delay/reception beam selection unit 207. Input signals and the user 1 signal to user L signal are multiplexed, and the multipath components of the user signals by propagation delays are also multiplexed.\nThe reception beam 1 path detection unit 2061 to reception beam M path detection unit 206M can also detect paths and measure the reception qualities of user signals at the detected path delays by using only a known symbol (pilot symbol or the like) of the user signal.\nThe path delay/reception beam selection unit 207 receives pieces of reception quality information of user signals corresponding to path delays/reception beam numbers as outputs from the reception beam 1 path detection unit 2061 to reception beam M path detection unit 206M. The path delay/reception beam selection unit 207 selects a path delay/reception beam number pair used for demodulation on the basis of the reception quality of the user signal, and outputs reception quality information of a user signal corresponding to the selected path delay/reception beam number to the transmission beam selection unit 209 and demodulation unit 210.\nThe transmission beam selection unit 209 receives the reception quality information of the user signal corresponding to the path delay/reception beam number as an output from the path delay/reception beam selection unit 207, and outputs the number of a transmission beam in the same direction as that of a reception beam having the delay path excellent in reception quality to the user 1 transmission beam switching circuit 2121.\nThe number of selected transmission beams is generally smaller than the number of path delay/reception beam number pairs used for demodulation. In many cases, the number of transmission beams is 1 in order to reduce interference with another user by transmission of a plurality of beams.\nThe demodulation unit 210 receives the reception quality information of the user signal corresponding to the path delay/reception beam number as an output from the path delay/reception beam selection unit 207, performs a demodulation process on the basis of the input path delay/reception beam number, and outputs user 1 reception data.\nThe user 1 modulation unit 2111 to user L modulation unit 211L respectively receive user 1 transmission data to user L transmission data (user transmission data), perform a modulation process, and output the modulated signals to the user 1 transmission beam switching circuit 2121 to user L transmission beam switching circuit 212L.\nThe user 1 transmission beam switching circuit 2121 to user L transmission beam switching circuit 212L receive the user 1 transmission beam number to user L transmission beam number as outputs from the transmission beam selection units 209 for respective users, and the modulated user signals as outputs from the user 1 modulation unit 2111 to user L modulation unit 211L. The user 1 transmission beam switching circuit 2121 to user L transmission beam switching circuit 212L select transmission beam formation units corresponding to the transmission beam numbers for the users from the transmission beam 1 formation unit 2131 to transmission beam J formation unit 213J, and output the modulated user signals to the selected transmission beam formation units.\nThe transmission beam 1 formation unit 2131 to transmission beam J formation unit 213J receive outputs from the user 1 transmission beam switching circuit 2121 to user L transmission beam switching circuit 212L, form fixed transmission beams different between the transmission beam 1 formation unit 2131 and transmission beam J formation unit 213J for the input signals, and output the fixed transmission beams to the antenna 1 radio transmission unit 2141 to antenna K radio transmission unit 214K. The number of fixed transmission beams, the shape of the fixed transmission beam, and the fixed transmission beam formation method are not specifically limited. An example of the shape of the fixed transmission beam is a quadrature multi-beam, and an example of the fixed transmission beam formation method is a method of multiplying input signals by a fixed complex beam weight by digital calculation (digital beam forming). In FIG. 7, the transmission beam 1 formation unit 2131 to reception beam J formation unit 213J are arranged on the input side of the antenna 1 radio transmission unit 2141 to antenna K radio transmission unit 214K, and form beams for digital signals of the baseband. A beam formation method in the radio band such as Butler matrix can also be adopted.\nEach of the antenna 1 radio transmission unit 2141 to antenna K radio transmission unit 214K comprises an amplifier, bandpass filter, mixer, local oscillator, quadrature modulation, low-pass filter, digital/analog converter, and the like. The antenna 1 radio transmission unit 2141 will be exemplified. The antenna 1 radio transmission unit 2141 receives outputs from the transmission beam 1 formation unit 2131 to transmission beam J formation unit 213J, performs reception processes such as digital/analog conversion of an input signal, quadrature modulation, frequency conversion from the baseband to the radio band, and amplification of a signal, and outputs the resultant signal to the transmission antenna element 2161.\nThe transmission array antenna 215 is formed from the K transmission antenna elements 2161 to 216K. Each of the transmission antenna elements 2161 to 216K is free from any limitation on the horizontal and vertical directivities, and has, for example, omnidirectivity or dipole. The K transmission antenna elements 2161 to 216K are arranged close to each other so as to correlate the transmission signals of the antenna elements with each other. The transmission array antenna 215 is not limited in the layout as far as the K reception antenna elements 2161 to 216K are arranged close to each other. An example of the layout is a circular layout or linear layout at the half-wavelength interval of the carrier.\nThe K transmission antenna elements 2161 to 216K receive and transmit signals in which user signals (user 1 signal to user L signal) by transmission beams as outputs from the antenna 1 radio transmission unit 2141 to antenna K radio transmission unit 214K are multiplexed.\nThe conventional multi-beam transmitter/receiver shown in FIG. 7 selects and receives a reception beam having a delay path excellent in reception quality from a plurality of fixed reception beams in reception. In transmission, the multi-beam transmitter/receiver selects and transmits a transmission beam in the same direction as that of a path delay/reception beam number pair excellent in reception quality from path delay/reception beam number pairs selected upon reception. With this process, the multi-beam transmitter/receiver can form a beam which increases the transmission/reception gain in a desired signal direction and decreases it in another direction.\nA problem of the conventional multi-beam antenna transmitter/receiver as shown in FIG. 7 is deterioration of the transmission characteristic. This is because a transmission beam in the same direction as that of a path delay/reception beam number pair excellent in reception quality is selected from path delay/reception beam number pairs selected upon reception, and a transmission beam optimum in the multipath environment cannot be selected. In the multipath environment, the user signal component contains a plurality of multipath components. These signal components generally arrive from different directions, and each reception beam contains a plurality of multipath components.\nThe conventional multi-beam antenna transmitter/receiver selects a transmission beam in the same direction as that of a path delay/reception beam number pair excellent in reception quality from path delay/reception beam number pairs selected upon reception. When the overall reception qualities of reception beams are compared, a reception beam different from a selected reception beam may exhibit a higher overall reception quality. The overall reception quality is prepared by calculating (e.g., adding) some or all of the reception qualities of multipath components (path delays) contained in a reception beam. An optimum transmission beam is a transmission beam in a direction in which the transmission beam coincides with (is identical to) or is close to a reception beam excellent in overall reception quality. The conventional multi-beam antenna transmitter/receiver cannot select any transmission beam optimum in the multipath environment.\nThis will be expatiated with reference to numerical values, but the present invention is not limited to these values.\nAssume that the path delay/reception beam selection unit 207 selects two upper pairs (pair a and pair b) from the following four path delay/reception beam number pairs.\nReception quality of pair a (path delay a/reception beam 1):10\nReception quality of pair b (path delay b/reception beam 2):8\nReception quality of pair c (path delay c/reception beam 2):5\nReception quality of pair d (path delay d/reception beam 1):1\nAt this time, if the transmission beam selection unit 209 selects one transmission beam, the conventional multi-beam antenna transmitter/receiver shown in FIG. 7 compares the reception qualities of pairs a and b (10>8), and selects a transmission beam in the same direction as reception beam 1. However, reception beam 2 has a higher overall reception quality obtained by calculating reception qualities for each reception beam (overall reception quality of reception beam 1=10+1< overall reception quality of reception beam 2=8+5). The conventional multi-beam antenna transmitter/receiver cannot select an actually optimum transmission beam."}
{"text": "It is common to use cryptography to protect or hide sensitive data found in software from being seen or used by third parties. For example, Information Rights Management (IRM) is a feature found in the Office® software package from Microsoft Corp. of Redmond, Wash. IRM can help prevent sensitive information from being distributed to or read by persons who do not have permission rights to access the content of the sensitive information. As another example, the Outlook® application, included in aforementioned Office® software package, enables users to create and send e-mail messages with restricted permission to help prevent messages from being forwarded, printed, or copied and pasted. Documents, workbooks, and presentations that are attached to messages with restricted permission are automatically restricted as well.\nSometimes, in addition to data that is found in software, the actual software product itself can be highly sensitive, such as, for example, in the case of software tools for Digital Rights Management (DRM) software systems. With the increased complexity of software products, it is becoming more difficult to control access to software products which may have many portions/sections, some of which are very sensitive to some users and some which may not be sensitive at all to other users."}
{"text": "This invention relates generally to the field of computer systems, and in particular, to sharing translation lookaside buffer (TLB) entries among multiple logical processors.\nComputing systems use a variety of techniques to improve performance and throughput. One technique is known in the art as multiprocessing. In multiprocessing, multiple processors perform tasks in parallel to increase throughput of the overall system.\nA variation of multiprocessing is known in the art as multithreading. In multithreading, multiple logical processors, which may comprise a single physical processor or multiple physical processors, perform tasks concurrently. These tasks may or may not cooperate with each other or share common data. Multithreading may be useful for increasing throughput by permitting useful work to be performed during otherwise latent periods, in which the performance level of the overall system might suffer.\nAnother technique to improve performance and throughput is known in the art as pipelining. A pipelined processor performs a portion of one small task or processor instruction in parallel with a portion of another small task or processor instruction. Since processor instructions commonly include similar sequences of component operations, pipelining has the effect of reducing the average duration required to complete an instruction by working on component operations of multiple instructions in parallel.\nOne such component operation is a translation from virtual addresses to physical addresses. This operation is often performed by using a translation lookaside buffer (TLB). It is a function of the TLB to permit access to high-speed storage devices, often referred to as caches, by quickly translating a virtual address from a task, software process or thread of execution into a physical storage address.\nIn systems which permit multiprocessing, including those systems that permit multithreading, identical virtual addresses from two different threads or software processes may translate into two different physical addresses. On the other hand, multiple threads or software processes may share a common address space, in which case some identical virtual addresses may translate into identical physical addresses. To prevent mistakes in accessing high-speed storage, the data may be stored according to physical addresses instead of virtual addresses.\nIf a high-speed storage device is accessed by multiple logical processors, the size of the TLB may be increased to allow storage of virtual address translations for each logical processor or thread of execution. Unfortunately, the time required to perform a virtual address translation increases with the size of the TLB, thereby reducing access speed and overall system performance. Alternatively, smaller faster TLBs may be physically duplicated for each logical processor, but physically duplicating these hardware structures may be expensive. Furthermore, in cases where multiple threads or software processes share a common address space, the TLB entries may include duplicates of some virtual address translations, thereby wasting space in this expensive resource. Providing private TLBs in a multithreaded processor, therefore, inefficiently uses this resource and prevents the logical processors from sharing translations when they share code or data. The inability to share translations is particularly harmful to the performance of multithreaded software, such as a database, wherein the logical processors often run threads that share a single address space. Sharing TLBs allows logical processors to dynamically partition the available resources based on the run-time needs of each processor and share translations, leading to more efficient use of the resource."}
{"text": "1. Field of the Invention\nThe present invention relates to game devices and more particularly a spherical dice apparatus wherein the random numbers which may turn up when a die is rolled are located on the outer surface of a sphere as opposed to the sides of a cube as is well known in a conventional die. The filed of the present invention further relates to modifications to a die to vary the method by which a random number may be turned up upon the roll of the spherical die.\n2. Description of the Prior Art\nThe conventional and well known die is a small cube marked on each face with from one to six spots and used, usually in pairs, in various games and in gambling by being shaken and thrown to come to rest at random on a flat surface.\nIn general, variations on the convention die concept have been created and patented in the prior art. The following patents and patent publications are illustrated of the known prior art:\n1. U.S. Pat. No. 973,595 issued to Wahlin in 1910 for \"Spherical Dice\". PA1 2. U.S. Pat. No. 4,546,978 issued to David in 1985 for \"Dice And Games\". PA1 3. U.S. Pat. No. 1,593,907 issued to Madan in 1925 for \"Game Device\". PA1 4. U.S. Pat. No. 4,807,883 issued to Silverman in 1989 for \"Game Apparatus And Dice Construction Therefor\". PA1 5. U.S. Pat. No. 4,164,351 issued to Baker in 1979 for \"Die-Resembling Game Cube\". PA1 6. U.S. Pat. No. 4,682,778 issued to Willis et al in 1987 for \"Political Game Utilizing Die With Interchangeable Faces\". PA1 7. U.S. Pat. No. 3,400,932 issued to Conrad in 1965 for \"Lawn Dice Having Finger Holes\". PA1 8. U.S. Pat. No. 809,293 issued to Friedenthal in 1906 for \"Game Apparatus\". PA1 9. German Offenlegungsschrift No. 3425 102 A1 published in 1986.\nU.S. Pat. No. 3,400,932 to Conrad for \"Lawn Dice Having Finger Holes\" illustrates a dice with a spherical member 10 disposed therein. The dice has a hollow interior having the shape of a truncated octahedron which contains a ball having an eccentric center of gravity. This is basically a lawn die which can be rolled by hand and which will automatically right itself on one of the flat outer faces. This spherical member has an eccentric weight embedded in the surface sot hat the ball rolling around within the hollow interior will impart a wobbly or eccentric motion to the die which will not only provide entertainment for use thereof, but will render it impossible for the player to predict by expertness of rolling same which of the surfaces will come up.\nU.S. Pat. No. 4,682,778 to Willis is relevant because it shows a changeable eight-sided die. The section inside the die are magnetically attached.\nU.S. Pat. No. 973,595 to Wahlin for \"Spherical Dice\" also discloses a spherical dice. A spherical body 5 is inserted inside the ball. The spherical body is made of heavy material. The sphere also contains various recesses inside it so that the ball can fall on one of the recesses thereby causing a number on the sphere to come up.\nU.S. Pat. No. 4,546,978 to David for \"Dice And Games\" discloses die configurations which have an outer shell in which a balanced weight is positioned. Also provided is a die construction with manual and/or chance adjustments.\nU.S. Pat. No. 1,593,907 to Madan for \"Game Device\" discloses a die construction of interest in as much as die body 6 is disposed within die 5. The interior die rolls within the exterior die and is also transparent so that both die faces can be seen.\nThe U.S. Pat. No. 4,807,883 to Silverman is relevant in that it illustrated various shapes of dice in FIGS. 5, 6, 7, and 8.\nU.S. Pat. No. 4,164,351 to Baker illustrated a dice assembly with a ball in it.\nU.S. Pat. No. 809,293 to Friedenthal for \"Game Apparatus\" discloses an odd shaped piece of dice.\nGerman Offenlegungsschrift discloses a dice having a hollow interior.\nWhile overall the concept of having an interior rotating device is shown as described in a number of the above Patents, what is not shown is the concept of having an interchangeable internal piece which causes the circular die to land in a particular location by having a circular weight land inside the interchangeable piece. What is further not shown is the concept of an interchangeable external shell.\nThere is a need for the improvements of the present invention which permit games of chance to be modified in difficulty through variation of the playing piece such as the dice."}
{"text": "In a conventional image-forming device, such as a laser printer, having a main device body, a process cartridge including a photosensitive drum and a developer unit is detachably mounted in the main device body. An access opening is formed in a prescribed surface of the main device body, while a cover is disposed on the main device body for exposing and covering the access opening. This cover is rotatably supported about a shaft extending along the prescribed surface of the main device body. By rotating the cover open about the shaft, the access opening is exposed, enabling a user to linearly remove a process car-ridge from or mount a process cartridge into the main device body through the access opening. When the cover is closed about the shaft, the cover covers the access opening.\nHowever, in an image-forming device having this construction, the cartridge cannot be mounted in or removed from the main device body unless the cover is opened wide. To open the cover wide, a large space is needed around the periphery of the main device body. Hence, the image-forming device occupies a greater space than the space required for installation, resulting in restrictions on installation locations.\nIn view of the foregoing, it is an object of the present invention to provide an image-forming device and a cartridge mounted in the image-forming device that are capable of Increasing the flexibility of choosing an installation location."}
{"text": "Many device testing applications, such as test applications for memory devices, require a test device capable of simultaneously testing many devices having multiple test channels. Such a test device typically requires periodic calibration to, among other items, calibrate various clock and data paths through the test device so that the test device delivers meaningful test results. Modem test devices are capable of testing many devices, and therefore have many, in some cases thousands, of test channels.\nA typical electronic circuit test device tests a device under test (DUT) by applying test signals of various logic states to input terminals on the DUT. The states of the signals that are produced at the output terminals of the DUT in response to the applied test signals are monitored to determine whether the DUT is behaving predictably.\nAn electronic circuit tester must maintain extremely accurate tolerances so that the performance of the DUT may be characterized accurately. Previous testing devices were manually calibrated using external measurement instrumentation. Other prior testing devices used a single calibration clock reference to calibrate all channels of the test device. Unfortunately, this requires that each channel of the test device be serially calibrated, thus requiring an undue amount of time to calibrate a test device having many hundreds or thousands of channels.\nAccordingly, a need exists for a way of reducing the amount of time required to calibrate a test device."}
{"text": "Wireless input devices such as the wireless mouse have become popular because they eliminate the inconvenience associated with wires and cables. Among their shortcomings is that a wireless mouse is heavy compared to a conventional wired mouse because it requires batteries to power up. Although initially the extra weight may not be bothersome, frequent users will feel the additional stress and strain on their hand and wrist. This can lead to, or aggravate, stress-related injuries such as carpal tunnel syndrome.\nReferring to FIG. 1A, there is shown a conventional computer system 100 using a wired mouse 102. The mouse derives its name from the resemblance to the animal of the same name. The wire is said to resemble the tail. Computer mice such as mouse 102 can detect the X-Y translation of the mouse position by mechanical or optical methods. The mouse movement is typically translated by a device driver (software that processes input from a device such as a mouse) to a relative movement of the cursor on the computer screen. As the mouse 102 is moved around a surface such as pad 104, the screen cursor moves as well in a well known relationship.\nComputer system 100 includes a processing system (computer 101) which communicates with the input/output (I/O) devices. Examples of conventional I/O devices include a keyboard, mouse, display, removable and fixed recordable media (e.g., floppy disk drives, tape drives, CD-ROM drives, DVD-RW drives, Flash drives, etc.), touch screen displays and the like.\nFIG. 1B illustrates the computer system of FIG. 1A with a typical wireless mouse 103. The wireless mouse 103 includes an optical system to detect X-Y translation of the mouse position. Mouse 103 also includes a radio transmitter to send wireless signals to the computer 101 when the mouse position has changed. These signals cause the computer 101 to update the location of a cursor displayed on a screen. Both the optical system and the radio transmitter are operated by batteries in the mouse 202. This mouse operates only on the pad 104, or on other surfaces that provide enough contrast so that the optical system can detect a change when the mouse is translated. If the mouse is lifted above the surface or placed on a highly reflective surface the optical system typically fails to detect movement in the X-Y plane correctly. The optical mouse 103 does not provide information on Z-translation (height) because conventional computer mice movement is limited to the x-y plane (e.g. movement over a mouse pad or a table).\nUS Published Application 2005/0219208 by Eichenberger et al. discloses a wireless input device, wherein a key stored in a receiver or a memory within the mouse can be transmitted by a radio-frequency identification (RFID) device provided in the mouse. This way, the mouse may be carried in a user's pocket and may be used to authenticate the user. However, the device disclosed uses conventional location determining means.\nUS Published Application 2005/0237295 by Anderson discloses a computer and RFID-based input device, wherein the input device can also operate together with another RFID product. For example, the computing arrangement can be adapted to receive and decode the signal from the RF transponders only when the computing arrangement detects the presence of an authorized RFID tag, which may be provided on a user ID card. However this device does not provide a radio frequency indication of the location of the mouse.\nUS Published Application 2006/0107061 by Holovacs discloses providing secure access to a KVM (keyboard, video, mouse) switch and other server management systems, wherein an RFID transceiver is coupled to or integrated with the system for detecting the presence of RFID tags that are within range of communications. A user workstation is only operable if the RFID transceiver detects a valid RFID tag.\nUS Published Application 2007/0195058 by Lin discloses a wireless cursor pointing device, which relates to a wireless cursor pointing device without the need for a battery. The wireless cursor pointing device transmits a signal through a Radio Frequency Identification (RFID) system. This device does not use RFID capability for location determining.\nU.S. Pat. No. 7,027,039 issued to Henty discloses a computer system with a passive wireless mouse, wherein the passive mouse includes passive transponder tags configured adjacent to the perimeter of the mouse ball comprising a body portion of the mouse. Further, an oriented antenna design, inductive coupling pattern or capacitor pattern are provided in tags chosen so that the tuning and activation of one tag is provided by the vertical motion of the mouse ball whereas the other tag will be affected by the horizontal ball motion. Therefore, as the mouse ball is rotated the two tags will be tuned and detuned repeatedly with a duty cycle which corresponds to the velocity of the rotating mouse ball. This thus provides two components of the rotational motion of the mouse ball which can be used to provide the specific speed and direction of the mouse ball. This in turn can be used to accurately track mouse ball position to provide conventional mouse control functions. However this system lacks the capability to detect movement in three dimensions.\nU.S. Pat. No. 7,280,097 to Chen et al. discloses a human interface input acceleration system, wherein in one implementation, the auxiliary input accelerator device may be a programmable radio frequency identification (RFID) key fob. Today, many workers gain access to their secure workplace via an RFID key fob. The circuitry in such a key fob may be housed inside an input accelerator with buttons, sensor arrays and a wireless communication protocol to be used as an auxiliary input accelerator device to the host device.\nTherefore, there is a need for a pointing device that overcomes the shortcomings of the existing art. More specifically, none of the devices mentioned above uses the radio frequency (RF) signal to detect translation of the device or to acquire the device position."}
{"text": "1. Field of the Invention\nThe present invention relates to a novel anionic polymerization for polymerizing a monomer in a solvent in the presence of a polymerization initiator. More particularly, it relates to an anionic polymerization for producing a polymer having an average molecular weight being substantially the same with the predetermined one, in high reproducibility by substantially removing a polymerization inhibitor in the polymerization system by a special manner before charging the monomer and the polymerization initiator.\n2. Description of the Prior Arts\nA living anionic polymerization is a significantly important polymerization process for producing a block or graft copolymer for a thermoplastic elastomer; an oligomer having both terminal functional groups being useful as a liquid elastomer or a polymer having uniform molecular weights useful as a macromolecular standard for determining molecular characteristics.\nIn an ideal living anionic polymerization in the absence of any side reaction beside the initiation and the propagation reaction, a molecular weight M of a polymer corresponds to a ratio of an amount of a charge monomer m(g) to an amount of a polymerization initiator I(mole). That is, the molecular weight of the polymer M is given by the equation (1): EQU M=m/I (1)\nHowever, in general, it has been difficult to remove completely a polymerization inhibitor which deactivates a polymerization initiator and an active living end from a solvent for polymerization, a monomer for polymerization and a polymerization atmosphere.\nIn such case, a molecular weight of the polymer M is given by the equation (2): EQU M=m/(1-.alpha.) (2)\nwherein .alpha. (mole) represents an amount of the polymerization inhibitor; which is a sum of .alpha.sa of an amount of the polymerization inhibitor remained in the solvent and the atmosphere and .alpha.m of an amount of the polymerization inhibitor remained in the monomer.\nWhen an anionic polymerization is carried out in a small scale in a laboratory, the solvent and the monomer used for the polymerization are small enough to attain a complete purification of the solvent and the monomer. A purification process is carried out in high vacuum so as to remove substantially the polymerization inhibitor remained in the solvent and the monomer by repeatedly contacting with a metal hydride hydrate, an alkali metal and an organometallic compound whereby the value .alpha. can be remarkably reduced. Therefore, in such case, it is possible to obtain a polymer having a molecular weight being substantially the same with the molecular weight given by the equation (1) so far as the molecular weight is not too high.\nOn the other hand, in a case of a bench scale or a factory scale of the polymerization, an amount of the solvent is too much to carry out in high vacuum and accordingly, the polymerization is carried out in an inert gas and it is impossible to purify completely the solvent as in high vacuum. The value .alpha. can not be negligible in comparison with an amount of the polymerization initiator I. Therefore, a molecular weight of the resulting polymer is higher than a predetermined molecular weight given by the equation (1).\nWhen a predetermined molecular weight of the object polymer is high, such as higher than 500,000, the molecular weight of the resulting polymer is remarkably higher than the predetermined molecular weight given by the equation (1). Sometimes, a polymerization does not proceed.\nIn a large scale process, it is difficult to maintain constant of the amount of the polymerization inhibitor .alpha.sa remained in the solvent and the atmosphere for polymerization in polymerization batches whereby it is difficult to expect a reproducibility on the relation of the predetermined molecular weight and the molecular weight of the resulting polymer."}
{"text": "Vehicles employ various power-sources for propulsion. Such power-sources may include an internal combustion engine and/or one or more electric motors or a fuel-cell.\nEach of such power-sources typically requires an energy storage device to receive and store an energy charge, and to supply the stored energy in order to sustain operation of that power-source. For example, a fillable tank may be employed for supplying fuel to run the internal combustion engine or the fuel-cell, and a chargeable electrical battery may be employed for supplying current to run an electric motor.\nA specific degree of the energy storage device's energy charge generally supplies the vehicle with a finite driving-range. Such a driving-range typically depends on a number of factors which may be related to the vehicle itself, as well as to road and weather conditions. A vehicle operator's driving style may also influence the vehicle's available driving-range on a specific degree of charge."}
{"text": "Conventionally, a clock recovery (clock reproduction) technique of recovering the clock signal from a received signal has been generally used in a receiving apparatus as a method of generating the clock signal required upon demodulation processing (for example, see Patent Document 1). FIG. 1 is a block diagram showing a main configuration of the receiving apparatus disclosed in Patent Document 1. In the receiving apparatus shown in FIG. 1, clock recovering circuit 13 recovers the clock signal from a received signal. The basic operation of clock recovery in clock recovering circuit 13 includes detecting points of level change of the received signal, controlling the frequency and phase of the clock signal held in clock recovering circuit 13 to synchronize with the detected points of level change and outputting the signal after control, as a recovered clock signal.\nPatent Document 1: Japanese Patent Application Laid-Open No. 2004-215220"}
{"text": "High Dynamic Range Imagery (HDRI) refers to processes, techniques, and digital image systems capable of reproducing real-world lighting and color data with high accuracy. HDRI initially was introduced to overcome dynamic range limitations in digital images, such as described in the scientific publication “Overcoming Gamut and Dynamic Range Limitations in Digital Images” (hereinafter “Ward98—2”), Proceedings of the Sixth Color Imaging Conference, November 1998, as well as in the scientific publication “High Dynamic Range Imagery” (hereinafter “Ward01”), Proceedings of the Ninth Color Imaging Conference, November 2001\nHDR pixel data commonly is represented using 96 bit data; 32 bits single precision IEEE floating point data for each RGB component. Standard digital images usually employ 24 bits per pixel; 8 bits for each RGB component. Thus, an HDR image has four times as much data as an LDR image. Therefore, due to relatively large amount of data for an HDR image, there is a need to substantially compress HDR image data.\nThe publication G. Ward, “Real Pixels,” Graphics Gems II, Ed. by J. Arvo, Academic Press, 1992 (hereinafter “Ward92”), describes the RGBE image format as representing the required dynamic range using 32 bits/pixel instead of 96 bits/pixel, providing a compression ratio of 3:1. Typical image and video compression algorithms, such as JPEG and MPEG, achieve much higher compression ratios, thus producing files hundreds of times smaller than the original source data. The RGBE format explicates a relevant compression problem by introducing the Exponent on a per-pixel basis, since even small errors that may be introduced by common compression methods, such as but not limited to JPEG and MPEG, generate exponentially higher levels of artifacts in the recovered image.\nEuropean Patent EP 0991020, owned by the Eastman Kodak Company, (hereinafter, “Kodak Patent”) describes a method to encode the extended color gamut of sRGB images by using residual images obtained by computing\nThe scientific publication “Perception-motivated High Dynamic Range Video Encoding”, Proc. of SIGGRAPH '04 (Special issue of ACM Transactions on Graphics), 2004 (hereinafter “Mantiuk04”) describes a method to encode HDR video into the MPEG4 standard. As described, the chrominance part of the image is expressed using a u′v′ color space similar to the color space used in the LogLuv image format discussed in Larson, G. W., “LogLuv encoding for full-gamut, high-dynamic range images,” Journal of Graphics Tools, 3(1):15-31 1998 (hereinafter “Ward98”), and also describes using 11 bits to store luminance data directly into the DCT coefficients. To overcome visible artifacts due to the augmented DCT range, a hybrid block coding system using additional data also is introduced\nThe scientific publication “High Dynamic Range Image and Video Data Compression”, IEEE CG&A, 2005 (hereinafter “Xu05”) describes a method to encode HDR image/video data, where, considering the RGBE image format, the RGB part of the image and the “E” (Exponent) part are separated and sent to different compression schemes. As described, the compression scheme for RGB is lossy while the scheme for “E” is lossless.\nThe publications “Subband Encoding of High Dynamic Range Imagery”, First Symposium on Applied Perception in Graphics and Visualization (APGV), August 2004, (hereinafter “Ward04”) and “JPEG-HDR: A Backwards-Compatible, High Dynamic Range Extension to JPEG”, Proceedings of the Thirteenth Color Imaging Conference, November 2005, (hereinafter “Ward05”) describe methods/systems for encoding HDR images and video in which a tone mapping operator is applied to the HDR source so that the HDR input is mapped smoothly into a 24 bit output domain with no components being clamped at 0 or 255. The original luminance of the HDR image is divided by the luminance of the tone-mapped LDR version, generating an HDR grayscale residual image which is log-encoded and used to restore the original frame.\nThe scientific publication “Backward Compatible High Dynamic Range MPEG Video Compression”, Proc. of SIGGRAPH '06 (Special issue of ACM Transactions on Graphics), 25 (3), pp. 713-723, 2006, (hereinafter “Mantiuk06”) discloses a method for compressing HDR video similar to that disclosed in Ward04 and Ward05.\nThe processes and techniques disclosed in each of the publications identified above disadvantageously call for additional coding and additional per-pixel operations, as well as additional pre- and/or post-correction processing, to achieve satisfactory results. Moreover, currently available MPEG and/or MPEG4 features and filters cannot be fully and/or easily exploited or employed by, or are incompatible with, many of the processes/systems mentioned above.\nConsumer digital displays today are only capable of displaying 8 bits for each RGB color component (i.e. 24-bit image data) and, considering luma/chroma separation, are only capable of displaying 8 bits for luma data that contain the high-frequency details of an image. With the introduction of high definition (HD) resolutions and bigger display screens—up to 100 inches—the limited precision of standard 24-bit digital images reveals several artifacts on such displays, even when uncompressed. By employing common compression solutions, like the JPEG or MPEG plethora of algorithms, these artifacts become even more visible to the human eye, due to the reduced precision, in number of bits, of the resulting decompressed output image.\nIn order to fully exploit the potential of HD resolutions, a display should be capable of showing 1920 shades, i.e. the maximum number of pixels in the longest pixel row at the highest HD resolution (1080p). With 8-bit data it is only possible to display 256 shades and moreover with MPEG compression this figure is usually greatly reduced, producing visible banding artifacts when using HD resolutions even on a 26 inches display. MPEG—and the latest H.264 standard—also introduces block-artifacts due to the nature of the algorithms which encode and store image data in blocks of 4×4 pixels up to 16×16 pixels, which are independent from each other, i.e. the output pixel luma value of a block has no correlation with an adjacent pixel of a different block. This usually leads to a high luminance contrast between edges of the surrounding blocks.\nIn order to avoid block-artifacts a number of algorithms commonly known as “deblocking filters” have been introduced. Such algorithms though are by nature very expensive in terms of the computing power required, which is linearly related to the dimensions of the image to be filtered. With HD resolutions, the number of pixels to be processed with a deblocking filter is up to 4.6 times the number of pixels being processed in a standard PAL or NTSC frame, thus a deblocking algorithm requires much more computing power in order to be applied to HD resolutions.\nHence, there is also a need to provide an effective method that enables encoding and decoding of higher precision image data (>24 bits/pixel) that is compatible with standard compression methods, such as JPEG and MPEG, and for a method that is able to remove typical image compression artifacts using the lowest possible computing power."}
{"text": "1. Field of the Invention\nThe invention relates to a door opening and closing device for opening and closing a door of, for example, a numerical control (NC) machine by using a moving robot.\n2. Description of Related Art\nIt is known that an NC machine is surrounded by a housing so as to prevent lubricating oil, chips, and similar materials from being scattered around during machining. The housing is provided with a large-sized sliding door so that the workpiece to be machined can be loaded into the housing and the machined workpiece removed from the housing using a robot.\nThe robot is an articulated robot, having a plurality of arms, and it is usually located in the vicinity of the door. A leading one of the arms is provided at its free end with a hand for gripping the workpiece. The workpiece to be machined is moved, by operating the arms, from the outside of the housing to a working position in the housing and the machined workpiece is moved again, by operating the arms, from the working position to the outside of the housing.\nIn moving the workpiece, using the arms of the articulated robot as mentioned above, it is necessary to open the sliding door to create an opening through which the workpiece and the arms can pass. A large-sized air cylinder is necessary for opening and closing the large-sized sliding door. If the articulated robot is operated to enter the housing under the condition where the sliding door is not sufficiently opened, there is a risk that the articulated robot will contact the sliding door. Accordingly, it is necessary to provide sensors for detecting opposite edges of the opening as defined by the open condition of the sliding door.\nAs mentioned above, in loading or unloading the workpiece into or from the housing, by operating the arms of the articulated robot, and machining the workpiece in the housing of the NC machine, it is necessary to open and close the large-sized sliding door and it is accordingly necessary to provide the large-sized air cylinder for driving the large-sized sliding door. Further, it is also necessary to provide a device for driving the large-sized air cylinder. In addition, it is necessary to provide at least two sensors for detecting the opposite edges of the opening defined by the open condition of the sliding door. For these reasons, the door opening and closing device in the prior art is complicated in construction and causes an increase in cost.\nIn addition, in such a device for opening and closing the sliding door by using an air cylinder, the operating speed of the sliding door and the positioning of the sliding door cannot be precisely controlled so that the period of time (cycle time) for loading and unloading of the workpiece by the articulated robot cannot be shortened."}
{"text": "Generally speaking, two types of rectifying schemes may be adopted in a secondary side of an isolated switching converter: (1) non-synchronous rectifying with a diode (FIG. 1A), and (2) synchronous rectifying with a synchronous rectifier, e.g., an N-MOSFET (FIG. 1B). The power dissipation-current characteristic is plotted in FIG. 2, for a diode (curve 12) and a synchronous rectifier (curve 11). In practical applications, the work area of a low power isolated switching converter always falls into the shadowed area. In the shadowed area, curve 11 is always above curve 12, i.e., the power dissipation of a diode is higher than the power dissipation of a synchronous rectifier. So, compared with a diode, a synchronous rectifier is more preferable because of less power waste and better efficiency in the low power isolated switching converter. Moreover, the temperature characteristic of the synchronous rectifier is better than that of the diode because of less power dissipation.\nSynchronous rectifiers have thus found increasingly wide applications in devices sensitive to power efficiency, such as laptop adapters, wireless equipment, LCD power management modules, power over Ethernet, and so on.\nIn the synchronous rectified switching converter, a voltage across the synchronous rectifier may be adopted to determine the on and off of the synchronous rectifier. However, the ringing of the voltage across the synchronous rectifier after the synchronous rectifier is turned off may cause mis-trigger."}
{"text": "Vertical hydroponic or aeroponic structures are known in the art.\nFor example, U.S. Pat. No. 6,021,602 issued to Orsi on Feb. 8, 2000 teaches a modular structure for aeroponic cultivation comprising a plurality of prefabricated vertical panels which have openings in order to receive cultures therein. U.S. Pat. No. 4,869,019 issued to Ehrlich on Sep. 26, 1989 teaches a self contained aeroponic apparatus having a reservoir for containing nutrient solution and a substantially vertical support. European patent application No. 0533939 filed on Mar. 5, 1992 teaches a perforated panel that has a substantially vertical orientation with a small inclination. The holes in the panel are configured to contain absorbing tampons for water in order to support plants. International application WO 02/09500 A1 teaches a hydroponic system having a vertical wall made of hollow tubes and an irrigation system as well as modules elements for receiving plants therein. European patent application No. 1210868 filed on Jan. 27, 2000 teaches a panel for supporting plants, the panels are mounted in between guard rails, their back sides provide for exposing the roots of the plants so as to be pulverized via spray nozzles with nutrient and water, their front sides are exposed to a light source."}
{"text": "1. Field of the Invention\nThe present invention relates to a component material for constructing a casing for various kinds of electronic devices.\n2. Description of the Related Art\nFor a casing for various kinds of electronic devices, such as a camera, a mobile phone, a PDA (Personal Digital Assistant), a personal computer, and so on, use has conventionally been made of materials such as metal (aluminum, stainless steel, titanium and magnesium) or a resin (acrylonitrile-butadiene-styrene (ABS) resin, polycarbonate, acrylic resin). These materials have been adopted not only from the functional consideration such as their moldability, rigidity and corrosion resistance, but also for aesthetics.\nRecently, further attention has also been paid to using, as a familiar, and environment-friendly material, wood or a woody material obtained by processed wood. For example, Jpn. Pat. Appln. KOKOKU Publication No. 7-24353 discloses an electromagnetic shielding material obtained through impregnation of wood or woody material with phenol resin and subsequent carbonization. The Publication discloses the idea of, while using the wood material, imparting an electromagnetic shielding property to a product obtained."}
{"text": "1. Field of the Invention\nThe present invention relates to a program recording system suitable for instructing from remote places such as a place away from home a recording and reproducing apparatus at home to record desired programs, using, for example, a communication terminal such as a mobile phone unit provided with a function of receiving a program of television broadcast (hereinafter, also called television) and the like, and relates to a communication terminal and a recording and reproducing apparatus which are used in the system.\n2. Description of the Related Art\nRecently, a communication terminal such as a mobile phone unit integrally provided with a portable television broadcast receiver (hereinafter called a television receiver) capable of receiving ground-based analogue television broadcast has been in practical use, and with the start of mobile broadcast and ground-based digital television broadcast, the life style in which television programs can be viewed and listened to anywhere and anytime has further become popular as a general trend.\nOn the other hand, a television program recording apparatus has been proposed, in which recording reservation of one or more programs can automatically be carried out by a user being at a remote place using a mobile phone unit to transmit one or more program discriminating data to a recording apparatus at home, when a television program to be recorded is known in advance (refer to Patent document 1, for example). According to this television program recording apparatus, the recording reservation of television programs can be performed from the remote place, which is very convenient.    Patent document 1: Japanese Published Patent Application No. 2001-285768\nHowever, the technology described in the patent document 1 is the one which performs recording reservation of a program that is decided to record in advance, so that the technology can not cope with the case where a user desires to record the program being viewed at present using the above mobile phone unit provided with a function of receiving television broadcast. Further, the technology may cope with the above case, when connecting a separate video recording apparatus to the above mobile phone unit capable of receiving television broadcast, when preparing a memory device of large capacity, or when carrying a recording apparatus and a memory device together; however, a recording apparatus must be connected every time and a recording apparatus and a memory device are required to be always carried with a user, which is not user-friendly and increases the cost.\nTherefore, a user who desires, after getting home, to again view the program that was viewed using a mobile phone unit with a television receiving function has not been satisfied with simple and inexpensive means. Further, there is nothing to meet the request of a user having viewed a program only halfway, who desires to view the program from the beginning or to view the remainder thereof at home."}
{"text": "The invention relates to a radial roller head according to the earlier part of patent claim 1.\nFrom DE 42 36 085 there is known a radial roller head with rollers arranged in a holder distributed at equal angular distances and at axial distances about the rolling axis. The roller surfaces have a spiral-shaped increasing course over the circumference opposed to their rotational direction. The rollers are coupled to a toothed gearing and a locking latches automatically after each complete roller revolution. Before the rolling procedure the locking is released, this being by the subject which is to be formed with the rollers. With the known radial roller head there is provided a stop which limits the introduction movement of the subject between the rollers, which is coupled to the locking of the rollers and which is displaceable in the direction of the rolling axis about an operating stroke against the action o)f compression springs. The locking is released on bearing of the subject due to its movement about the operating stroke.\nWith the known radial roller head there is the danger that during the releasing procedure or also thereafter a relative movement takes place between the subject and rollers if the subject has not achieved its stationary situation. This procedure naturally inhibits the forming procedure."}
{"text": "Silicon-containing films are used for a wide variety of applications in the semiconductor industry. Silicon-containing films include silicon films such as polycrystalline silicon (poly-Si) and epitaxial silicon, silicon germanium (SiGe), silicon germanium carbide (SiGeC), silicon carbide (SiC), and silicon nitride (SiN). As circuit geometries shrink to ever smaller feature sizes, lower deposition temperatures are preferred, for example because of introduction of new materials into semiconductor devices and reduction of thermal budgets of shallow implants in source and drain regions. Moreover, it is evident that non-selective (blanket) and selective deposition of silicon-containing films will be needed for future devices. For example, semiconductor fabrication requires tight specification limits on thickness and resistivity for epitaxial silicon films. Epitaxial silicon deposition can be a first step in a process flow where the crystal lattice of the bulk silicon is extended through growth of a new silicon-containing layer that may have a different doping level than the bulk. Matching target epitaxial film\nthickness and resisitivity parameters is important for the subsequent fabrication of properly functioning devices.\nOne example of the use of selective deposition of epitaxial silicon-containing films is for manufacturing silicon-on-insulator (SOI) devices with raised source and drains. During SOI device fabrication, processing may consume an entire silicon film in source and drain regions, thereby requiring extra silicon in these regions that can be provided by selective epitaxial growth (SEG) of silicon-containing films. Selective epitaxial deposition of silicon-containing films can reduce the number of photolithography and etching steps that are needed, which can reduce the overall cost and complexity involved in manufacturing a device. Despite the preference for lower temperature deposition processes, thermal deposition of epitaxial silicon using the traditional silane (SiH4) and dichlorosilane (DCS, SiCl2H2) source gases generally require high deposition temperatures (e.g., greater than about 850-900° C.) to achieve deposition rates that are high enough for the process to be incorporated into processes for manufacturing of devices. Moreover, the traditional silane and dichlorosilane source gas processes have limited deposition selectivity with respect to different substrate materials. Thus, the present inventors have recognized that improved methods are needed for deposition of silicon-containing films onto substrates that permit selective deposition and deposition at reduced temperatures."}
{"text": "1. Field of the Invention\nThe present invention relates to reaction products of polyisobutenes having an average degree of polymerization P of from 10 to 100 and a content E of from 60 to 90% of double bonds which can be reacted with maleic anhydride, where E=100% would correspond to the theoretical value for the case where each molecule of the polyisobutene has such a reactive double bond, with oxides of nitrogen or mixtures of oxides of nitrogen and oxygen.\nThe present invention furthermore relates to the use of these reaction products as fuel and lubricant additives and to gasoline engine fuels and lubricants containing these reaction products.\n2. Discussion of the Background\nThe carburetor and intake system of gasoline engines, as well as injection systems for metering fuel into gasoline and diesel engines, are contaminated with impurities which are caused by dust particles from the air, uncombusted hydrocarbon residues from the combustion space and the crankcase vent gases passed into the carburetor.\nThe residues shift the air/fuel ratio in the idling state and in the lower part-load range so that the mixture becomes richer and the combustion more incomplete and in turn the proportions of uncombusted or partially combusted hydrocarbons in the exhaust gas increase and the gasoline consumption rises.\nIt is known that these disadvantages can be avoided by using fuel additives for keeping valves and carburetor or injection systems clean (M. Rossenbeck in Katalysatoren, Tenside, Mineraloladditive, Editors J. Falbe and U. Hasserodt, page 223 et seq., G. Thieme Verlag, Stuttgart, 1978).\nDepending on the mode of action, but also on the preferred place of action of such detergent additives, a distinction is now made between two generations of such assistants.\nThe first generation of additives could prevent only the formation of deposits in the intake system but could not remove existing deposits, whereas the additives of the second generation can accomplish both (keep-clean and clean-up effect) and can do so because of their excellent heat stability, in particular in zones at relatively high temperatures, ie. in the intake valves.\nIn general terms, the molecular structural principle of fuel detergents can be expressed as the linking of polar structures with generally relatively high molecular weight, nonpolar or lipophilic radicals.\nTypical members of the second generation of additives are often products based on polyisobutenes in the nonpolar moiety. Among these in turn, additives of the polyisobutylamine type are particularly noteworthy.\nU.S. Pat. No. 3,576,742 (1), filed in 1968 and granted in 1971, describes reaction products of branched long-chain aliphatic olefins, for example polypropylene, polyisobutylene or copolymers of ethylene and isobutylene, and oxides of nitrogen as detergents for lubricants. These olefins are prepared from lower olefins of 2 to 6 carbon atoms by conventional polymerization methods. In the resulting nitro-containing reaction products, the presence of further functional groups, such as hydroxyl, nitroso, nitrate, nitrite or carbonyl, and the relative proportions thereof are unknown. The structure and composition of these olefins are thus undefined.\nDE-C 27 02 604 (2) published in 1978 describes a preparation process for polyisobutenes having an average degree of polymerization of from 10 to 100, by means of which highly reactive polyisobutenes, ie. those having predominantly terminal double bonds, were obtained for the first time. On the other hand, the polyisobutenes prepared by conventional polymerization methods, as in (1), have only a low content of terminal double bonds.\nThe nitro-containing reaction products disclosed in (1) and based on conventional polyisobutene do have a certain action as lubricant additives, but this action is still unsatisfactory. Furthermore, it is said that the compositions can also be used as additives for fuels for gasoline engines."}
{"text": "Printing apparatuses such as printers, copying machines, and facsimile apparatuses are designed to print an image formed from a dot pattern on a printing material such as a paper sheet or thin plastic plate on the basis of image information. Such printing apparatuses can be classified into inkjet printing apparatuses, wire dot printing apparatuses, thermal printing apparatuses, laser beam printing apparatuses, and the like in accordance with their printing schemes. Of these schemes, the inkjet scheme (inkjet printing apparatus) is designed to print an image by discharging ink (printing liquid) droplets from the orifices of a printhead and causing the droplets to fly and land on a printing material.\nAs many printing apparatuses are used in recent years, they are required to realize high-speed printing, high resolution, high image quality, and low noise level. An inkjet printing apparatus can meet these requirements. In an inkjet printing apparatus which prints an image by discharging ink from a printhead, the temperature of ink at the discharge section has a strong influence on stable ink discharge and stable ink discharge amount, which are necessary for meeting the above requirements. More specifically, if the ink temperature is too low, the ink viscosity excessively increases. This makes it impossible to discharge the ink by normal discharge energy. Conversely, if the ink temperature is too high, the discharge amount excessively increases to cause ink overflow on a printing paper sheet, resulting in poor image quality.\nTo prevent this, conventional inkjet printing apparatuses have a temperature sensor at the printhead portion and employ a method of controlling the ink temperature at the discharge section within a desired range on the basis of the detected temperature of the printhead or a method of controlling a discharge recovery process.\nAs a heater for temperature control, a heater member joined to the printhead portion is used. Some inkjet printing apparatuses which print an image by forming flying droplets using thermal energy, i.e., some apparatuses which discharge ink droplets by growing bubbles by film boiling of ink, use a discharge heater itself as a heater for temperature control. Especially, when the discharge heater is used, the heater must be energized not to grow ink bubbles.\nIn a printing apparatus which obtains ink droplets to be discharged by forming bubbles in solid ink or liquid ink using thermal energy, the discharge characteristic largely changes depending on the temperature of the printhead. It is therefore particularly important to manage the ink temperature at the discharge section and the printhead temperature that greatly influences the ink temperature.\nIn managing the printhead temperature, it is important to measure the ink temperature at the discharge section which has a large influence on the discharge characteristic. However, it is very difficult to measure the ink temperature because the discharge section also generates heat, and this makes the detection temperature of the temperature sensor largely vary more than the variation in ink temperature necessary for management, and also because the ink itself moves. For these reasons, even when a temperature sensor is simply arranged near the printhead to accurately measure the ink temperature at the time of discharging, the variation in temperature of the ink itself can hardly be measured.\nAn inkjet printing apparatus has been proposed, which stabilizes the ink temperature indirectly by stabilizing the printhead temperature as an ink temperature management means. U.S. Pat. No. 4,910,528 discloses an inkjet printer which has a means for stabilizing the printhead temperature by predicting a subsequent discharge heater driving amount in a predetermined time using, as a reference, the detection temperature of a temperature sensor arranged near the discharge heater.\nMore specifically, the apparatus stabilizes the printhead temperature by controlling, in accordance with the predicted temperature, a printhead heating means, a discharge heater energization means, a carriage driving control means for maintaining the printhead temperature to a predetermined value or less, a carriage scanning delay means, a means for reducing the carriage scanning speed, a means for changing the printing sequence of ink droplet discharge from the printhead, and the like.\nHowever, the inkjet printer disclosed in U.S. Pat. No. 4,910,528 has a measurement error of the temperature sensor and a time lag between reading by the temperature sensor and reflection of the reading result on driving. It cannot sufficiently stabilize discharge in recent high-speed printing. In addition, since the printhead incorporates the temperature sensor, the cost of printhead increases. Hence, no inexpensive printing apparatus can be provided."}
{"text": "The present invention relates to a disposable diaper having a fastening system for connecting a front waist and a rear waist of the diaper to each other.\nJapanese Laid-Open Utility Model Application No. Hei7-22725 discloses a disposable diaper having tape fasteners extending outward from transversely opposite side edges of a rear region of the diaper adapted to be releasably fastened to receiving sheets bonded to an outer surface of a front region of the diaper, wherein graduations are marked on the receiving sheets in circumferentially around a waist of the diaper. These graduations serve as a positioning indicator when the tape fasteners are fastened to the receiving sheets so that these graduations may be utilized to maintain a tightening force around a wearer's waist at a constant level.\nThe receiving sheets are generally made of a plastic film or nonwoven fabric and the graduations are printed on surfaces of these receiving sheets. Particularly in the case of a baby diaper, materials used for the diaper should be carefully selected from the viewpoint of a sanitary security, since a baby often places their fingers or various other objects into their mouth. While hygienically innocuous ink is used for printing the graduations, it is desired to avoid contact between the ink and a baby's mouth.\nAccordingly, it is a principal object of the invention to eliminate a demand of printing the graduation and to prevent regions occupied by the receiving sheets from becoming unacceptably rigid by utilizing a pattern in which an adhesive for bonding the receiving sheets to the front region is applied as the graduations."}
{"text": "1. Field of the Invention\nThis invention relates to jewelry boxes and, more specifically, to jewelry boxes of the type that contain a hidden compartment or compartments therein which would be relatively easy to access, and yet to young children or a burglar it would not be known that there are hidden compartments within the jewelry box.\n2. Description of the Related Art\nJewelry boxes of one type or another have been around for ages and have taken various forms but there has always been a need to have a jewelry box or the like that would have an area which would be imperceptible and wherein objects could be placed and be hidden from view so as not to be too easily or readily accessible to young children or those with larceny in their heart.\nAs far as is known there is no prior art that teaches a jewelry box that would have an open front to receive at least one slidably positionable drawer wherein the drawer would not fill the entire interior content or chamber of the jewelry box to thereby allow the positioning of one or more compartments within the interior of the jewelry box which would not be readily discernable to an outsider or young child, but yet would be readily accessible to the owner of the jewelry box.\nIn order for a hidden compartment or compartments to be located within a jewelry box and yet still be readily accessible, the hidden compartment of necessity must have at least some side walls that would be congruently configured to the compartment within which they are contained and wherein the other features of the jewelry box such as a drawer would be of lesser width in dimension than the overall width of the jewelry box so as to permit the positioning of a hidden compartment within the interior of the jewelry box and wherein the interiorly received container may be pivotally mounted for ease of positioning wholly within the container or rotatably or pivotally movable to an accessible location partially outside the interior of the jewelry box.\nA search of related art has been conducted and no specific reference has been found that would detract from the overall patentability of the herein disclosed jewelry box having at least one hidden compartment which may or may not be pivotally mounted for positioning wholly within the interior of the jewelry box and then would be positionable partially exteriorly for ease of access.\nU.S. Pat. Nos. 2,738,250 and 4,195,727 each disclose jewelry boxes having a pair of outwardly pivoting trays.\nU.S. Pat. Nos. 2,564,122 and 5,823,328 each disclose storage boxes such as for jewelry having pivoting compartments.\nU.S. Pat. No. 6,322,123 and more specifically the element 24 disclosed in FIGS. 2-5 inclusive, disclose a storage drawer having pivoting trays.\nU.S. Pat. Nos. 5,127,719 and 6,059,388 each disclose jewelry boxes having concealment features, but none of the type as disclosed in the herein disclosed invention.\nSimilarly, U.S. Design Pat. Nos. 203,305, 270,980, 451,671, 454,687 and 455,902 fail to disclose anything that would be remotely similar in over all concept with respect to hidden compartments to the instantly disclosed invention."}
{"text": "1. Field of the Invention\nThis invention relates to novel pharmaceutical compositions for the improvement of dysuria such as pollaki-uria or incontinence of urine. Dysuria is often caused by adult diseases such as cerebrovascular disorders and prostatic diseases, which results in not only serious hamper of the patients' action in their daily life but also impairment of the patients' dignity. On the other hand, it gives a great burden on the nursing family. As the world becomes more aged in population, and also from the viewpoint of quality of life, development of agents for the improvement of dysuria with more definite efficacy and higher safety is highly desired. The compositions according to the invention are useful as a therapeutic agent with higher safety in dysuria such as pollakiuria and incontinence of urine.\n2. Description of the Prior Art\nPharmaceutical compositions for the treatment of pollakiuria and incontinence of urine in urology include cholinolytics, antimuscarinics and calcium antagonists. None of such composition may be satisfactory either in efficacy or in adverse reaction. For example, whereas, cholinolytics typical of which is butylscopolamine are observed to be inhibitory on cystic smooth muscles, their use is associated with disadvantages in that they are also highly active on other cholinergic viscera such as digestive tracts, cardiac muscles, sialaden and pupillary muscles. Since treatment for dysuria needs a long period of time for continued use of a drug even in aged patients, it should not only be therapeutically highly effective but also be highly safe."}
{"text": "1. Field\nExample embodiments relate to a transmissive electrode, a photoelectric device and/or an image sensor including the same.\n2. Description of the Related Art\nA photoelectric device may convert light into an electrical signal. Examples of photoelectric devices include a photodiode, a phototransistor, and the like. A photoelectric device may be applied to an image sensor, a solar cell, an organic light emitting device, and the like.\nAn image sensor including a photodiode may have higher resolution and accordingly a smaller pixel. At present, a silicon photodiode may be used. Silicon photodiodes may have less sensitivity when the pixel is smaller due to a smaller absorption area. Accordingly, an organic material that is capable of replacing silicon has been researched.\nOrganic materials may have a high extinction coefficient and may selectively absorb light in a particular wavelength region depending on a molecular structure, and thus may simultaneously replace a photodiode and a color filter and improve sensitivity and contribute to high integration.\nIn order to realize a photodiode having a higher photoelectric conversion efficiency based on the organic material, it is desirable to improve characteristics of a transmissive electrode."}
{"text": "1. Field of the Invention\nThe invention pertains to a system for producing the natural aerobic biological decomposition of human waste by controlling the distribution of air and liquid within a conversion tank where waste is subjected to aerobic bacteria, moisture and oxygen.\n2. Description of the Related Art\nIn the installation of restrooms in parks and recreational areas, the disposal of human waste is often a serious problem. It may be necessary to locate the restroom where electric power is not available, or the terrain and subsoil may be such as to prevent the installation of septic tanks or more conventional waste disposal systems. Restrooms in mountain areas built upon rock cannot utilize a septic tank system, and the disposal of the waste by direct discharge may violate environmental regulations. \"Dry\" toilets and similar waste conversion systems are known, and some systems utilize chemicals and other catalysts for breaking down the waste and facilitating decomposition. However, chemical dry toilets require considerable attention and maintenance, and are not suitable for use in remote locations, such as recreational parks and hiking areas. As environmental regulations become increasingly stringent, the need for waterless on-site waste conversion systems capable of meeting environmental standards increases, and the advantages of natural aerobic biological decomposition systems using dehydration and evaporation becomes important.\nIt is knownto devise \"dry\" toilets by atmospheric venting, such as shown in U.S. Pat. No. 3,136,608, and it is also known to use a plurality of evaporating tanks with a waste conversion system such as shown in U.S. Pat. No. 3,624,665. However, such disclosures have a number of drawbacks, and require considerable maintenance and attention to properly function.\nForced hot air is commonly used in dry toilets to augment evaporation and decomposition, as shown in U.S. Pat. Nos. 4,196,477 and 4,313,234, but such units require energy to heat air and are not suitable for remote installations."}
{"text": "This invention relates to methods for the preparation of mixtures of long chain monounsaturated alcohols, such as jojoba alcohol. This oily liquid and its compositions are used for the topical treatment of subdermal infections such as herpes virus infections, and the transdermal delivery of pharmacological agents for the treatment of various disorders.\nA large percentage of the world population are infected with herpes viruses. Three of the most common herpes viruses are herpes simplex virus-1 (HSV-1) which is the cause of facial and ocular sores, herpes simplex virus-2 (HSV-2) which has a predilection for genital areas, and herpes simplex virus-3 (HSV-3), also named herpes zoster or varicella zoster, which causes chicken pox and later shingles. Once an individual is infected, herpes viruses become latent principally in nerve cells, and can reactivate to cause recurrences of the original symptoms. Herpes migration to the brain or spinal cord leads to encephalitis and meningitis, which are life-threatening conditions. Human herpes virus-8 (HHV-8) is associated with the skin cancer Kaposi sarcoma. These herpes episodes are each susceptible to topical treatments because the viruses replicate in subdermal cells during a recurrence, and before an eruption into a lesion. There are several treatment options for herpes infections but no cures. Several of the nucleoside analog drugs can be effective if taken prophylactically on a daily basis. They are less effective if administered at the time of the recurrence, either orally or topically. The nucleoside drugs inhibit viral replication by penetrating into the cell and interfering with nucleic acid production. They are not virucidal, and depend on a functional immune system to deactivate any virus present. New anti herpes virus agents, their targets and therapeutic potential have been reviewed (Alrabiah and Sacks, Drugs, 52: 17-32, 1996), as have topical treatments of herpes simplex virus infections (Hamuy and Berman, Europ. J. Dermatol., 8:310-319, 1998; Evans and Tyring, Dermatol. Clinic, 16: 409-419, 1998; Syed et al., Clin. Drug Invest., 16: 187-191, 1998).\nAlcohols with chain lengths of 16 to 20 carbon atoms and 1 to 4 double bonds inhibited herpes simplex and another lipid enveloped viral bacteriophage in cell cultures (Sands et al., Antimicrob. Agents Chemother., 15: 67-73, 1979). These unsaturated alcohols were more potent in vitro than saturated alcohols with shorter chain lengths (Snipes et al., Antimicrob. Agents Chemother., 11: 98-104, 1977). A patent (Rivici et al., U.S. Pat. No. 4,513,008, 1985) describes the inhibition of enveloped viruses such as herpes with linear polyunsaturated acids, aldehydes or primary alcohols with chain lengths of 20 to 24 carbons and 5 to 7 double bonds. These reports were followed by the investigations and development of n-docosanol as a topical treatment for herpes infections\nn-Docosanol, also named 1-docosanol or behenyl alcohol, is a straight chain 22 carbon saturated alcohol, which occurs in the bark, flowers and fruit of the tree Pygeum africanum. n-Docosanol is reported to have broad spectrum in vitro activity against lipid enveloped viruses such as herpes (Katz)et al., Proc. Nat. Acad. Sci., 88:10825-10829, 1991; Katz et al., Ann. N.Y.Acad. Sci., 724: 472-488, 1994; Pope et al., J. Lipid Res., 37: 2167-2178, 1996; Pope et al., Antiviral Res., 40:85-94, 1998), and also the human inmmunodeficiency virus HIV (Marcelletti et al., AIDS Research and Human Retroviruses, 12: 71-74, 1996). These studies demonstrated that the antiviral activity of n-docosanol includes inhibition of fusion with or viral entry into the cell, while being mediated by intracellular metabolic biotransformation of the drug. A series of patents on the composition of mixtures of n-docosanol in formulations that render it useful for topical application supports these published reports (Katz, U.S. Pat. No. 4,874,794, 1989; Katz, U.S. Pat. No. 5,071,879, 1991; Katz, U.S. Pat. No. 5,166,219, 1992; Katz, U.S. Pat. No. 5,194,451, 1993; Katz, U.S. Pat. No. 5,534,554, 1996). n-Docosanol is not virucidal but is virustatic, interfereing with viral replication, and dependent on a functional immune system to destroy viruses. n-Docosanol is a crystalline waxy solid insoluble in water which needs to be formulated with a surfactant and carrier to facilitate dermal penetration and interaction at the target cell level, This limitation was also noted where several other virustatic long chain compounds with 18 plus linear carbons including amides, alkanes, acids and alcohols needed to be formulated with a surfactant and carrier to facilitate penetration of the epidermis (Katz et al., U.S. Pat. No. 5,534,554, 1996; Katz et al., PCT WO098/11887, 1998). The latter two reports claim a composition of n-docosanol or other long chain compounds with a surfactant and a pharmaceutically acceptable diluent or carrier as the active viral replication inhibitor, rather than the pure individual compounds. The preparation of n-docosanol is not reported in these patents. The long chain alcohols and other compounds are generally waxy solids, and would not be expected to penetrate skin layers alone without a carrier. In a study using 10% n-docosanol suspended in an aqueous system containing a non-ionic surfactant and a carrier, mean healing time of lesions in humans infected with herpes labialis (HSV-1) was shortened (Habbema et al., Acta Derm. Venereol., 76: 479-481, 1996). A 12% n-docosanol cream was tested as a possible transmision prophylactic of simian immunodefficiency virus (SIV) in rhesus macque monkeys (Miller et al., Antiviral Res., 26: A277, 1995). Intravaginal application before exposure prevented transmission in five of the six animals tested. n-Docosanol and other saturated alcohols with chain lengths of 20 to 26 carbons reportedly promote corneal healing due to eye injury (Muller, U.S. Pat. No. 5,214,071, 1993; Muller, U.S. Pat. No. 5,296,514, 1994).\nJojoba oil has been available commercially for more than twenty years, and several million pounds are used in cosmetic formulations annually. Jojoba oil is a mixture of esters composed principally of both long chain mono unsaturated alcohols and carboxylic acids (Wisniak, The Chemistry and Technology of Jojoba Oil, publ. by American Oil Chemists Society, Champaign, Ill., 272 pp, 1987). A significant characteristic of jojoba oil is its ability to be absorbed quickly by the skin. Extensive testing and use of jojoba oil has established that it is completely non-toxic when applied to human sin, or administered orally to mice, rats, marmots and rabbits (Taguchi and Kunimoto, Cosmetics and Toiletries, 92: 53-61, 1977; Clark and Yermanos, Biochem. Biophys. Res. Commun., 102: 1409, 1981; Hamm, J. Food Sci., 49: 417-428, 1984; Verschuren and Nugteren, Food Chem. Toxicol., 27: 45-48, 1989). Humans who have ingested jojoba seeds, which are 50% oil, have not been harmed, although some nausea occurred when as much as 200 grams were eaten. In mice, jojoba oil has functioned as an intestinal lubricant (Verbiscar et al., J.Agric. Food Chem., 28: 571-578, 1980). It is estimated that about 20% of jojoba oil is split by esterases in the gastrointestinal system, thus producing jojoba alcohol in situ. After dermal absorption, jojoba oil would be at least partially metabolized to jojoba alcohol. Jojoba oil is a generally recognized as safe for cosmetic uses (Final Report on the Safety Assessment of Jojoba Oil and Jojoba Wax, J. Amer. Coll. Toxicol., 11(1): 57-74, 1992).\nJojoba alcohol has been prepared from jojoba oil by hydrogenolysis with sodium and alcohol (Molaison et al, J. Amer. Oil Chem. Soc., 36: 379-382, 1959). In this reaction, the carboxylic acid part of the ester is converted to its corresponding alcohol, in comparison with chemical hydrolysis where the fatty acids remain intact and must be separated from the alcohols in the mixture. Hydrogenolysis doubles the amount of jojoba alcohol that can be obtained from jojoba oil. One jojoba alcohol product prepared by hydrogenolysis resulted in a mixture of 6% octadec-9-enol, 62% eiocos-11-enol, 28% docos-13-enol and 4% tetracos-15-enol (Taguchi, Proc. Sixth Int. Conf. Jojoba and Its Uses, eds. Wisniak and Zabicki, Ben-Gurian Univ. Negev, Beer-Shiva, Israel, p 371-391, 1984). The actual alcohol composition will vary according to the source of jojoba oil. The relative amounts of individual alcohol components in jojoba alcohol depend on the ester composition of jojoba oil, a product obtained from seeds of the desert shrub Simmonsia chinensis harvested in the Southwestern United States, Mexico, Israel and South America.\nPure long straight chain monounsaturated alcohols are sometimes waxy liquid or low melting solids, but when in a mixture as in jojoba alcohol exist as a colorles odorless oil at normal ambient temperatures. A characteristic of jojoba alcohol is that it is readily absorbed by human skin, much like the parent jojoba oil. Jojoba alcohol does not require a carrier or surfactant to facilitate transdermal penetration. Jojoba alcohol is reported in one patent on lipstick components along with a large number and variety of other compounds with diverse structures and functions seeking lipocolor properties which would make the compositions feel less dry (Sato, Lipocolor Composition, U.S. 5,700,453, 1997). Jojoba alcohol is mentioned in the specifications of one patent on kojic acid where it is an excipient (Honda, U.S. Pat. No. 5,750,563, 1998).\nThe safety of jojoba alcohol for cosmetic uses has been reported (Taguchi, Proc. Sixth Int. Conf. Jojoba and Its Uses, eds. Wisniak and Zabicki, Ben-Gurian Univ. Negev, Israel, p 371-391, 1984). Mouse, rabbit, marmot and human tests were made for jojoba alcohol confirming that this product is very safe for topical application. Mutagenicity tests were also negative.\nWhen applied to an incipient herpes simplex virus recurrence, jojoba alcohol quickly penetrates the epidermal layer to the subdermal cells where viral replication leading to symptomatic disease would otherwise occur (Verbiscar, PCT WO 9916245, Dec. 9, 1999, patents pending). Jojoba alcohol appears to function in the same manner that n-docosanol functions, that is by inhibiting lipid enveloped viruses from fusion with and entering into cells. Irritation is relieved and viral replication is delayed, while the host immune system is alerted to destroy the free virus units.\nWhen treated early in the prodrome or even the erythema inflammation stages, herpes blisters do not form or at least are inhibited in persons with functional immune systems. This inhibitory action is applicable to enveloped viruses which cause lesions at epidermal surfaces. Herpes simplex viruses which cause recurrent facial sores (HSV-1) and genital sores (HSV-2), shingles (HSV-3) and Kaposi sarcoma (HHV-8) are treatment targets.\nThe metal hydride process for producing jojoba alcohol (II) from jojoba oil (I) is an improvement over the reported sodium hydrogenolysis because yields are higher and the reaction is safer and easier to carry out. The key product prepared is a mixture described as monounsaturated alcohols containing 14 to 24 carbon atoms represented in the reaction: \nwhere m and n are 5 to 13, and the carbonxe2x80x94carbon double bonds are cis or trans. Jojoba alcohol is representitive of this general formula (II). The actual individual alcohol composition in the mixture can vary, as does the composition of esters in jojoba oil from which jojoba alcohol is derived. It is a characteristic of secondary plant metabolites such as the parent jojoba oil to vary in composition according to plant variety, pollination, soil and climate conditions of cultivation. Jojoba alcohol remains an oil at ambient temperatures above about 13 degrees centigrade because it is a mixture of monounsaturated alcohols. The melting point, or freezing point, is lowered by virtue of this being a mixture as is the melting point of any chemical compound which contains an impurity. Jojoba alcohol is not a single pure compound, and its freezing point behaves accordingly, maintaining it as an oil at ambient temperatures which are well below normal skin temperatures. This oily characteristic facilitates its applications and uses as a topical transdermal treatment product, without the need for a cofactor carrier or surfactant. The fact that jojoba alcohol is a mixture of individual pure compounds, and thus remains an oil at normal temperatures is significant to its functions. Mixtures of similar long chain monounsaturated alcohols can be prepared by hydrogenolysis of sperm whale oil, which is itself a mixture of mono esters of long chain monounsaturated carboxylic acids and long chain monounsaturared alcohols similar to jojoba oil. It is understood that mixtures of monounsaturated alcohols can be prepared by combining pure compounds similar to those in jojoba alcohol, and that these formulated mixtures will be oils and have properties similar to natural jojoba alcohol.\nIn the previously reported sodium hydrogenolysis, a coreagent alcohol and the jojoba oil ester are added to liquid sodium in boiling toluene. The system generates nacent hydrogen which reduces the ester groups to alcohol groups without reducing the carbonxe2x80x94carbon double bonds. In our hands, this process yielded a 73% yield of pure distilled jojoba alcohol. However, workup and purification of the final product was difficult due to the presence of carboxylic acid by-products resulting from partial hydrolysis of the ester rather rather than hydrogenolysis. The sodium salts of these long chain fatty acids act as surfactants and cannot be water-washed out of the mixture. In addition, while this reaction can be done readily in small quantities, large scale production with liquid sodium presents difficulties. To circumvent these problems, we invented a new chemical hydrogenolysis of the ester group of jojoba oil to the alcohols using a commercial reagent named sodium bis-(2-methoxyethoxy)aluminum hydride (Bazant et al., Tetrahedron Letters, 29; 3303-3306, 1968). With this reagent, a 90.5% yield of jojoba alcohol was obtained which is identical to the liquid sodium hydrogenolysis product, but with no dangerous reactants or workup problems. This jojoba alcohol freezing point was about 13 degrees centigrade, and remains an oil at normal ambient temperatures. Major alcohol components of jojoba alcohol are octadec-9-enol, eicos-11-enol, docos-13-enol and tetracos-15-enol. The actual amount of any individual alcohol in a jojoba alcohol preparation will vary according to the source of jojoba oil used in the hydrogenolysis. Other alcohols, ranging in chain length from 14-24 carbon atoms, may also be present as minor components (Wisniak, Chemistry and Technology of Jojoba Oil, publ. American Oil Chemists Society, Champaign, Ill., 272 pp, 1987). This process is also applicable to sperm whale oil, producing a similar mixture of long chain monounsaturated alcohols. Although the sodium bis(2-methoxyethoxy)aluminum hydride is available commercially, it can be prepared in situ for direct hydrogenolysis of jojoba oil. A convenient preparation reacts one mole of sodium plus one mole of aluminum with two moles of 2-methoxyethanol in benzene or toluene at temperatures above 100xc2x0 C. under hydrogen pressure (J. Vit, B. Cosensky and J. Machacek, French patent, 1,515,582). The disappearance of the insoluble sodium and aluminum metals indicates this reaction is complete and the reagent is formed. Jojoba oil can then be added slowly to the reagent solution, at low temperatures to complete the hydrogenolysis to jojoba alcohol. Workup of the final product is the same as for the metal hydride system described here.\nThe generic structure of the major components of this jojoba alcohol mixture and examples of individual mono unsaturated alcohols composing jojoba alcohol follows, where the double bond can exist in cis and trans forms and m and n can vary from 5 to 13 carbons.\nJojoba alcohol is a mixture of principally monounsaturated linear alcohols with 14 to 24 carbon atoms in the chain. A mixture of these alcohols is readily prepared by hydrogenolysis of jojoba oil, which is a mixture of esters. Mixtures of monounsaturated linear alcohols can also be prepared from other sources. Sperm whale oil is an example of mixtures of long chain esters of monounsaturated alcohols and carboxylic acids which can also be converted to the corresponding mixture of long chain monounsaturated alcohols using this metal hydride hydrogenolysis process. A mixture of the individual alcohol components can also be prepared by combining specific amounts of each alcohol. The more alcohols in the mixture, the more likely will the mixture be liquid at ambient temperatures. A formulated mixture of individual alcohols will act like jojoba alcohol as a transdermal penetration carrier and viral fusion inhibitor. Jojoba alcohol represents mixtures of these alcohols which will remain liquified at ambient temperatures, thus facilitating its use in topical transdermal treatments. While the specifications contain particular details, these should not be construed as limitations of the scope of the invention, but rather examples of preferred embodiments. Jojoba alcohol can actually be comprised of any percent variation of the individual long chain monounsaturated alcohols, depending on the source of the jojoba oil from which it is derived. There are many varieties of Simmonsia chinensis from which the jojoba oil is obtained from the seed, with different pollination source plants, and different cultural conditions. All of these differences contribute to the final chemical composition of the jojoba oil esters. The most significant structural factors in the jojoba alcohol compounds is that they are highly lipophylic with a hydrophylic hydroxyl end group, and a carbon carbon double bond in the inner chain which facilitates transdermal penetration."}
{"text": "There are many electricity transmission systems and distribution systems operating at a fixed frequency, for example, at 50 Hz or 60 Hz. Grid frequency control systems control resources such as power sources (generators) and energy storage resources.\nIn order to maintain grid frequency stability, grid operators send commands for controlling the amount of power withdrawn from the grid by loads and the amount of power injected to the grid by generators, storage, and renewable energy sources. The power sources and the energy storage resources control inputs or outputs of their power.\nHowever, large delay occurs in transmission of command signals from a grid operator's computer to receivers. To overcome the delay, there is a system that creates, at a receiver side, a control signal having a grid frequency, instead of using commands sent from the grid operator (see PTL 1, for example)."}
{"text": "Solid state lasers developed rapidly during the year since the first ruby laser in the 1950's. (T. H. Maiman--Nature 187, 493 (1960)) In 1964 the Nd doped YAG was discovered (J. E. Geusic, H. M. Marcos and L. G. Van Uitert (App. Phys. Lett. 4, 182 [1964]). This Nd-YAG technology is in wide use today and is being continuously improved. By the 1970's solid state lasers entered a period of mass commercialization. Today, many new solid state lasers are based on rare earth and other elements such as uranium, chromium, samarium, ytterbium, thulium, holmium, erbium, and many others are being developed (A. A. Kaminskii \"Crystalline Lasers: Physical Processes and operating Schemes\", CRC Press 1996).\nMany lasers are used by aircraft, tanks, vehicles as well as foot soldiers. Most of the these lasers are either pumped by flash lamps or laser diodes which are bulky. These processes are inefficient and expensive as well. Therefore, there is a need for a low cost, lighter weight, efficient laser. This patent describes lasers which use a thermally stimulated photon emitting light pipe as the pumping device, which may be heated by any heat source such as combustion, nuclear, radioisotopic, chemical, solar or the like.\nA variety of different combustion systems useful as a heat source for the thermally stimulated photon emitting light pipe or \"light pipe\" can be designed depending on specific application. Some combustion systems can employ an emissive matrix combustor having the beneficial characteristics of producing a high power density and providing a high radiant output heat source. Emissive matrix combustion systems provide an extremely high power density (up to 30 W/cm.sup.2) and have a radiant output of about 50% of total chemical energy that is released in a combustion process.\nAnother approach can be to design a vortex combustion system, having high power density and extended residence time of hot combustion products, in a light pipe containing section that increases the thermal energy transfer from the exhaust to the light pipe.\nIt is also possible to fabricate a microchannel combustor having a high power specific density (up to 100 W/cm.sup.3) and having a highly extended light pipe area, or combination of these or other combustion technologies such as open flame combustors, surface combustors, and the like.\nSuperemissive insulation that covers internal shell of the combustion chamber is an additional source of the narrow band photons which can transfer energy through the light pipe to the laser. Superemissive materials can be described as materials that, when heated to a threshold temperature, includes one or more electrons that jump to a different electron energy level in quantum increments which causes the emission of ultraviolet, visible or infrared radiation in a wavelength band related to the atom's inner electron shell vacancy. Emitted radiation produced as a result of such electron transition is often within a narrow band and can, therefore, be absorbed efficiently by a laser system, such as Nd doped YVO.sub.4 or YAG.\nThermally-stimulated superemissive materials produce radiation in relatively concentrated narrow special bands compared to blackbody or greybody emitters which typically exhibit a broadband thermal emission such as from a flash lamp. As a result of the concentrated, narrow spectral band, photons emitted from the superemitter and focused to the laser system one would expect greater efficiency than that generated by a blackbody emitter operating at the same heat flux. However, a blackbody emitter may be constructed to emit visible or infrared radiation in a narrow spectral band by using one or more band-pass filters interposed between the superemitter and the target laser or other device. Photon filtering, however, is not energy efficient. On the other hand, laser diode can emit in a very narrow band, but must be powered by power conditioned electric systems which are expensive, inefficient, heavy and bulky.\nPhoton generators that use superemitters to emit radiation and, therefore, generate photons are well-known and are disclosed in U.S. Pat. Nos. 4,776,895, 4,793,799, 4,906,178. For example, such a photon generator includes a porous ceramic burner in the shape of a cylinder having an annular passage extending therethrough. A superemitting fiber layer is disposed along the inside wall of the ceramic burner and is made from a high temperature fiber or coating comprising, for example, pure or doped oxides of uranium, thorium, ytterbium, aluminum, gallium, yttrium, erbium, holmium, zirconium, chromium or other high temperature oxides. When subjected to thermal energy, the fiber layer emits radiation that is directed to a central axis running along the annular passage of the cylindrical burner.\nAs the ceramic burner is heated, the superemissive fiber layer emits radiation that can be filtered, by use of a cylindrical hollow filter disposed within the ceramic burner, to emit radiation at a selected bandwidth. The radiation passing through the filter is directed to the surface of an optical cable that is disposed centrally within the annular passage of the ceramic burner and filter. The photons generated by the fiber layer are, therefore, directed onto the optical cable and are channeled through the cable to each cable end, which is directed to a target comprising a photovoltaic cell. Accordingly, the photons generated by the superemissive fiber layer within the ceramic burner are directed through the optical cable into the photovoltaic cell, and converted to electricity.\nThe embodiments of the photon generator disclosed in earlier identified U.S. patents and U.S. patent applications do not promote the most efficient generation and collection of photons to pump and power laser systems. In the above-discussed embodiments, the emission of photons is effected through application of thermal energy to the body, which passes through the body of thermal conduction to the superemissive material. Accordingly, a large amount of thermal energy is wasted through the mechanism of thermal conduction. Additionally, the photons emitted by the superemissive material are not collected efficiently, as relatively weaker photons, or photons that are emitted from a distance further away from the optical waveguide than other emitted photons, are not collected within the waveguide. Accordingly, many of the photons generated by the superemitter are not collected by the optical waveguide and, therefore, are not directed to the target(s).\nIt is, therefore, desirable that a device be constructed to facilitate the economic generation and collection of selected photons, and be capable of directing them to a laser to optically pump the laser. It is desired that the device be simple to construct and be easily adaptable for use with a number of different thermal generation sources. It is desirable that the device be configured to accommodate use with one or more target laser systems and also to generate electricity by conversion of photons to electricity by means of one or more photovoltaic cells. It is desirable that the device be constructed in such a manner as to eliminate or minimize the effects of thermal shock on the cell. It is desirable that the device be capable of being manufactured from conventional materials using current manufacturing techniques. It is also desirable for this unit to be low cost, light weight and compact as well as energy efficient. The ability to operate directly from the available heat source without the need for electric power and power conditioning is very desirable."}
{"text": "This invention relates to the disproportionation (metathesis) of olefins. In accordance with one aspect, this invention relates to a catalyst suitable for use in the disproportionation of olefinic hydrocarbons. In accordance with another aspect, this invention relates to a process for the disproportionation of olefinic hydrocarbons. In accordance with a further aspect, this invention relates to a catalyst suitable for use in the disproportionation of olefins comprising at least one of molybdenum oxide and tungsten oxide, an inorganic refractory oxide support, and at least one methylating agent. In accordance with a further aspect, this invention relates to a catalyst suitable for use in the disproportionation of olefins comprising a support and at least one of molybdenum oxide and tungsten oxide promoted with at least one methylating agent. In accordance with another aspect, this invention relates to a process for the disproportionation of olefinic hydrocarbons with a disproportionation catalyst as hereinbefore described under conditions of temperature and pressure which effect disproportionation of the feed.\nThe disproportionation, or metathesis, of olefins is a reaction in which one or more olefinic compounds are transformed into other olefins of different molecular weights. The disproportionation of an olefin with itself to produce an olefin of a higher molecular weight and an olefin of a lower molecular weight can also be referred to as a self-disproportionation. For example, propylene can be disproportionated to ethylene and cis-, and trans-2-butene. Another type of disproportionation involves the cross-disproportionation of two different olefins to form still other olefins. An example would be the reaction of one molecule of 2-butene with one molecule of 3-hexene to produce two molecules of 2-pentene.\nBy the term \"disproportionation\" or \"metathesis\" throughout this specification is meant the conversion of the feed olefinic (or unsaturated) hydrocarbon to a mixture of olefinic (or unsaturated) hydrocarbons having different numbers of carbon atoms than the feed hydrocarbons.\nAmong the catalysts that have been developed for disproportionation are those comprising inorganic refractory oxides containing a catalytic amount of at least one of molybdenum oxide and tungsten oxide. The present invention is based upon the discovery of a way to improve the activity of such a catalyst.\nAccordingly, an object of this invention is to provide a method for the conversion of olefins.\nAnother object of this invention is to provide a catalyst for the conversion of olefins.\nStill another object of this invention is to provide a method for converting olefins to olefins having different numbers of carbon atoms than the feed hydrocarbons.\nStill another object is to provide a method for improving the activity of a disproportionation catalyst for the conversion of olefins into olefins having different numbers of carbon atoms than the feed hydrocarbons.\nOther aspects, objects and the several advantages of the invention will be apparent to one skilled in the art upon reading the disclosure including a detailed description of the invention and the appended claims."}
{"text": "Differential mechanisms are conventionally used to transmit driving forces to a pair of drive wheels of a vehicle and to permit said wheels to change their speed of rotation and travel when the vehicle is turning. Under certain road conditions such as mud, ice, snow, or rough irregular terrain, the traction of one drive wheel with the ground surface may differ from the traction of the other drive wheel. Such difference in traction with the differential mechanism in operation may result in a loss of driving force imparted to one of the drive wheels. In some instances, one drive wheel may slip to such an extent that it may spin freely with a complete loss of driving force at that wheel. Substantial driving power is thereby lost. Under such conditions, locking of the differential to prevent differential operation permits transmission of driving forces through the locked differential mechanism directly and substantially equally to the wheels.\nLocking devices for differential mechanisms have been employed heretofore to alleviate such conditions. A simple prior proposed locking device has included a lock pin moved into engagement with a hole in one of the parts of the differential mechanism. The shear stress imparted to such a pin was very great and undesirable because such stress often resulted in shearing of the lock pin and a resumption of differential operation.\nAnother type of prior proposed locking device for a differential mechanism included the movement of a locking collar splined to a half axle and movable axially to engage a toothed locking member fixed for rotation with a rotatable differential cage, U.S. Pat. No. 3,215,000. Another prior proposed locking device for a tractor transmission included the axial movement of a sliding second pinion gear along a drive shaft having a fixed pinion gear for engagement with juxtaposed bevel gears on spaced brake drums, U.S. Pat. No. 2,836,084."}
{"text": "Hand-held labelers have been used to print bar codes as well as human readable information on a label that is applied to an article to designate price, department number, etc. When more than one piece of information, such as a plurality of bar codes, are to be entered into known labelers, data entry has been accomplished via a keyboard or a host computer. Bar codes typically include an information portion of the code plus guard codes. The guard codes are disposed at the start and end of the information portion of the bar code as well as in the middle thereof where the middle guard code identifies the direction that the bar code should be read. For keyboard data entry, an alphanumeric representation of the information portion of the bar code is entered via the keyboard and a digital representation thereof is stored in a list in the labeler's memory. In the case of data entry via a host computer, a digital representation of the information portion of the bar code is downloaded from the host computer and stored in a list in the labeler's memory. To select information such as a particular bar code from the stored list, a user may step through the list by actuating a key on the keyboard of the labeler while the labeler displays human readable representations of the code information stored in the list. Once code information is selected for printing the labeler assembles a digital representation of a bar code by adding guard codes to the code information. After assembling the code, the labeler may print it on one or any number of labels. Data entry via a keyboard of a large group of codes can be very tedious, time consuming and unreliable. Although data entry via a host computer is generally reliable and efficient, it is not always possible. Hand-held labelers are often used in stores where a user is required to enter information already printed on a shelf label into the labeler so that the same information can be printed on articles to be placed on the shelf. In this type of situation data entry via a host computer is generally not possible.\nHand-held labelers are also known to include an optical reader or scanner such as a light pen for use in inputting data to be printed on a label. U.S. Pat. No. 4,652,317 shows such a hand-held labeler wherein the light pen is used to scan a code. The code may be printed on a label secured to a shelf that is intended to support articles to be labeled with information corresponding to the shelf label. Data representing the scanned code is stored in a RAM and remains there until a subsequent label is scanned. After a single label is scanned, one or more labels may be printed with the scanned information. Although the task of data entry is considerably eased by the scanning operation of this labeler, as opposed to those limited to manual data entry via a keyboard, the labeler is still limited in that only one label may be scanned and stored at a time."}
{"text": "This invention relates to a small size multi-way waterproof connector in which a waterproof seal between a connector housing and a plurality of wires, each connected to a terminal, is enhanced.\nA related waterproof connector, which is disclosed in Japanese Patent Publication No. 11-329571A, will now be described specifically with reference to FIGS. 6 to 8. A connector housing 11 of the waterproof connector 10 includes an inner housing 12, which is integrally molded of synthetic resin, and has a plurality of terminal chambers 13, an outer housing 17 of synthetic resin, in which the inner housing 12 is fitted, and a spacer 28 of synthetic resin which is provided between the inner housing 12 and the outer housing 17 to retain female terminals 14\nAs shown in FIGS. 6 and 7, the inner housing 12 has a box-like body 12a, and upper and lower sides of this box-like body 12a are open at a rear portion thereof. The terminal chambers 13 are formed by a central horizontal wall 12b and upper and lower vertical side walls (serving as partition walls) 12c. The female terminals 14 ale adapted to be received in the terminal chambers 13, respectively. Retaining pawls 15 are integrally formed respectively on opposite side portions and a central portion of each of the upper and lower faces of the body 12a. Insertion holes 12d of a rectangular shape are formed respectively through those portions of a front wall of the body 12a opposed respectively to the terminal chambers 13, and male terminals in a mating connector (not shown) are passed through these insertion holes 12d, respectively. Two pairs of press connecting blades 14b are formed by stamping on opposite side walls of a box-like portion 14a of the female terminal 14 at a rear portion thereof.\nAs shown in FIGS. 6 and 7, the outer housing 17 includes an inner wall portion 17a of a substantially square tubular shape, an outer wall portion 17b of a substantially square tubular shape formed around the inner wall portion 17a, and a bottom wall portion 17c interconnecting rear end portions of the inner and outer wall portions 17a and 17b. Thus, the outer housing 17 has a double-wall box 2 shaped construction having an open front side. A central portion of the bottom wall portion 17c is thickened, and rubber plug chambers or recesses 19 of a circular cross-section With a large diameter are formed respectively in those portions of an inner face of this thickened portion opposed respectively to the terminal chambers 13, and waterproof rubber plugs 18 are inserted respectively into these rubber plug chambers 19 by press-fitting or the like. Wire insertion holes 21 of a circular cross-section with a smaller diameter are formed in the outer face of the thickened portion, and communicate with the rubber plug chambers 19, respectively, and wires 20 are passed through the wire insertion holes 21, respectively. The waterproof rubber plug 18 has corrugated, cylindrical inner and outer peripheral faces, and the wire 20 is passed through the bore of the rubber plug 18, with no gap formed therebetween.\nRetaining hole of a rectangular shape are formed respectively in opposite side portions o a front portion of each of upper and lower walls of the inner wall portion 17a of the outer housing 17, and the retaining pawls 15, formed respectively on the opposite side portions of the upper and lower faces of the body 12a of the inner housing 12, are releasably engaged in these retaining holes 22, respectively. A rectangular longitudinal retaining hole 23 is formed in a central portion of the front portion of each of the upper and lower walls of the inner wall portion 17a, and the retaining pawls 15, formed respectively on the central portions of the upper and lower faces of the body 12a of the inner housing 12, are releasably engaged in these retaining holes 23, respectively. A V-shaped packing-receiving part 25 for receiving an annular waterproof packing 24 of rubber is formed integrally on the outer face of the inner wall portion 17a of the outer housing 17 at the inner end portion thereof. Those portions of the inner face of the inner wall portion 17a of the outer housing 17, disposed immediately adjacent respectively to front edges of the retaining holes 22 and 23, are formed respectively into tapered faces 26. Retaining holes 27 are formed respectively in upper and lower walls the outer wall portion 17b of the outer housing 17, and elastic retaining arms of the mating connector (not shown) are releasably engaged in these retaining holes 27, respectively.\nAs shown in FIGS. 6 and 7, the spacer 28 includes a body 28a of a substantially square tubular shape for fitting in the inner wall portion 17a of the outer housing 17, a flange portion 28b of a substantially square tubular shape, integrally formed on and extending in a stepped manner rearwardly from a front end of the body 28a so a to fit on the inner wall portion 17a of the outer housing 17, and a bottom wall portion 28c closing a rear end of the body 28a. Thus, the spacer 28 has a box-like shape having an open front end.\nThe body 12a of the inner housing 12 is fitted into the tubular body 28a of the spacer 28. Rib-like projections 29 for preventing disengagement of the terminals are formed integrally on inner faces of upper and lower walls of the body 28a of the spacer 28, and these projections 29 serve to retain the rear edge of the body 12a of the inner housing 12, and also serve to retain the rear edges of the box-like portions 14a of the female terminals 14, received respectively in the terminal chambers 13, respectively.\nNotches 30 are formed respectively in those portions of an interconnecting portion (which interconnects the body 28a and the flange portion 28b of the spacer 28), corresponding respectively to the retaining pawls 15 on the inner housing 12, and also notches 31 are formed respectively in those portion of this interconnecting portion corresponding respectively to flanges 16 formed on the inner housing 12. A retaining pawl 32 is formed integrally on the outer face of each of the upper and lower walls of the body 28a of the spacer 28, and is disposed between the upper or lower) notches 30, and these retaining pawls 32 are releasably engages respectively in the retaining holes 23 in the outer housing 17. When the spacer 28 is completely fitted into the outer housing 17, the distal end of the flange portion 28a of the spacer 28 serves to hold the packing 24 retained on the packing-receiving part 25 on the inner wall portion 17a of the outer housing 17.\nWire insertion holes 33 are formed respectively through those portions of the bottom wall portion 28c of the spacer 28 corresponding respectively to the wire insertion holes 21 in the outer housing 17. When the spacer 28 is completely fitted into outer housing 17, the bottom wall portion 28c of the spacer 28 serves to hold the rubber plugs 18 inserted respectively in the rubber plug chambers 19 in the bottom wall portion 17c of the outer housing 17. As shown in FIG. 6, the wires 20 are passed respectively through the wire insertion holes 21 in the outer housing 17, the rubber plugs 18 and the wire insertion holes 33 in the spacer 28, a d each wire 20 is press-fitted into a space between each pair of press-connecting blades 14b of the female terminal 14 received in the associated terminal chamber 13, and the terminal chambers 13 and the wires 20 are sealed by the rubber plugs 18 and the packing 24.\nFor assembling the waterproof connector 10 of the above construction, the rubber plugs 18 are inserted respectively into the rubber plug chambers 19, formed in the inner face of the bottom wall portion 17c of the outer housing 17 (which forms the outer portion of the connector housing 11), and also the packing 24 is fitted in the packing-receiving part 25 on the inner wall portion 17a of the outer housing 17, as shown in FIG. 8A.\nThen, the body 28a of the spacer 28 is fitted into the inner wall portion 17a of the outer housing 17 as shown in FIG. 8B, so that the retaining pawls 32 on the body 28a of the spacer 28 are retaining engaged respectively in the retaining holes 23 in the inner wall portion 17a of the outer housing 17. As a result of the retaining engagement of the retaining pawls 32 (formed on the body 28a of the spacer 28) in he respective retaining holes 23 formed in the inner wall portion 17a of the outer housing 17, the rubber plugs 18 are retained against withdrawal by the bottom wall portion 17c of the outer housing 17, and also the packing 24 is retained against withdrawal by the slanting distal end of the flange portion 28b of the spacer 28, and therefore the waterproof ability of the connector housing is further enhanced.\nThen, the wires,, 20 are passed through the respective wire insertion holes 21 in the bottom wall portion 17c of the outer housing 17 from the outer side thereof, the respective rubber plugs 18 and the respective wire insertion holes 33 in the bottom wall portion 28c of the spacer 28, as shown in FIG. 8C. Then, each wire 20 is press-connected to the pairs of press-connecting blades 14b of the female terminal 14 received in the associated terminal chamber 13 in the inner housing 12 (which forms the inner portion of the connector housing 11), as shown in FIG. 8D.\nThen, the inner housing 12 is fitted into the body 28a of the spacer 28 as shown in FIG. 8E, so that the retaining pawls 15 on the body 12a of the inner housing 12 are retaining engaged respectively in the retaining holes 23 in the inner wall portion 17a of the outer housing 17, thus completing the assembling operation of the waterproof connector 10. At this time, the rear edge of the body 12a of the inner housing 12, as well as the rear edges of the box-like portions 14aof the female terminals 14, received respectively in the terminal chambers 13, is locked by the projection 29 formed on the inner faces of the upper and lower walls of the body 28a of the spacer 28, and the projections 29 will not be deformed outwardly, and therefore the disengagement of the female terminals 14 from the respective terminal chambers 13 is positively prevented. And besides,. since the retainment of the female terminals 14 and the retainment of the rubber plugs 18 are simultaneously achieved by the spacer 28, it is not necessary to provide any member used exclusively for preventing the withdrawal of the rubber plugs, and therefore the number of the component parts is reduced, thereby reducing the cost. Furthermore, the spacer 28 is retained in a double manner by the retaining engagement of the retaining pawls 32 of the spacer 28 in the respective retaining holds 23 in the outer housing 17 and by the retaining engagement of the retaining pawls 15 of the inner housing 12 in the respective retaining holes 23 in the outer housing 17, and therefore the disengagement of the rubber plugs 18 and the packing 24 is positively prevented, thereby further enhancing the reliability of the waterproof effect.\nIn the above related waterproof connector 10, however, the rubber plugs 18 are first received respectively in the rubber plug chambers 19 in the outer housing 17 as shown in FIG. 8B, and therefore when passing the wires 20 through the respective wire insertion holes 21 in the outer housing 17 and the respective wire insertion holes 33 in the spacer 28 as shown in FIGS. 8B and 8C and also when fitting the inner housing 12 into the spacer 28, fitted in the outer housing 17 while sliding the wires 21 relative to the wire insertion holes 33 in the spacer 28 and the wire insertion holes 21 in the outer housing 17 as shown in FIGS. 8D and 8E, the wires 20 are liable to be buckled because of the sliding friction between the wires 20 and the rubber plugs 18 as shown in FIG. 9, and this would sometimes affect the production of wire harnesses. Namely, during the wire passing operation for passing the wires 20 and during the housing-fitting operation for fitting the inner housing 12 into the outer housing 17, the operation load due to the sliding Friction between the rubber plugs 18 and the wires 20 increases, which leads to a possibility that the efficiency of the assembling operation of the waterproof connector 10 is lowered.\nAnd besides, the rubber plugs 18 are completely covered with the bottom wall portion 17c of the outer housing 17 and a peripheral wall portion 17d around the rubber plug chambers 19, and therefore during the assembling operation of the waterproof connector 10, it was difficult to confirm from the outer side of the outer housing 17 whether or not the rubber plug 18 was received in each rubber plug chamber 19.\nThe present invention has been made in order to solve the above problems, and an object of the invention is to provide a small-size multi-way waterproof connector, as well as a method of assembling the same, in which the efficiency of an assembling operation is enhanced.\nIn order to achieve the above objects, according to the present invention, there is provided a waterproof connector comprising:\na connector housing including a terminal chamber for accommodating a terminal to which an electric wire is connected;\na rear holder attached to a rear end portion of the connector housing;\na waterproof rubber plug having a through hole through which the electric wire is inserted, one end of which is configured as a cylindrical sealing portion, the other end of which is configured as a proximal end portion integrally provided with a rear end wall of the rear holder; and\na plug chamber formed in the rear end portion of the connector housing so as to communicate with the terminal chamber for accommodating the rubber plug therein when the rear holder is attached with the connector housing,\nwherein the dimension of the plug chamber is determined such that the sealing portion of the rubber plug is accommodated while being compressed by an inner wall of the plug chamber during the attaching operation of the rear holder.\nPreferably, a depth dimension of the plug chamber is determined such that a distal end of the sealing portion of the rubber plug is compressed in the attaching direction of the rear holder when the rear holder is attached with the connector housing.\nSince the rubber plug is formed integrally with the rear holder, in the assembling operation in which the wire is passed through the insertion hole of the rubber plug, and then the female terminal is connected to the end portion of the wire, and is received in the terminal chamber in the connector housing, and then the rear holder is fitted into the plug chamber of the connector housing, the insertion resistance bet ween the electric wire and the rubber plug can be efficiently decreased. Thus, the electric wire is prevented from being buckled, thus enhancing the efficiency of the assembling operation of the waterproof connector. Therefore, a wire harness, comprising a bundle of wires, can be easily produced.\nAnd besides, since the sealing portion of the rubber plug is accommodated while being compressed by an inner wall of the plug chamber during the attaching operation of the rear holder, a failure to attach the rubber plug is eliminated, and a, seal between the rubber plug and the plug chamber in the connector housing, as well as a seal between the rubber plug and the wire, can be formed easily and positively, thus enhancing the sealing performance.\nPreferably, the dimension of the plug chamber is so determined as to satisfy the following relationship:\nA less than Cxe2x88x92B+D\nwhere A is an inner diameter of the plug chamber; B is an inner diameter of the through hole of the rubber plug; C is an outer diameter of the sealing portion of the rubber plug; and D is an outer diameter of the electric wire.\nUnder the above condition, the efficiency of the assembling operation of the waterproof connecter, as well as the waterproof performance thereof, is further enhanced.\nPreferably, the connector housing includes a plurality of terminal chambers and the rear holder includes a plurality of waterproof rubber plugs associated with the terminal chambers respectively. The respective proximal ends of the rubber plugs are exposed to an outer face of the rear end wall of the rear holder and connected with each other via a groove formed on the outer face of the rear holder.\nAccordingly, the plurality of rubber plugs can be firmly secured to the rear end wall of the rear holder, and the rubber plugs will not be disengaged from the rear holder even when passing the wires respectively through the insertion holes of the rubber plugs. Therefore, the multi-way waterproof connector, having the excellent sealing performance, can be easily assembled in a short. time, and the efficiency of the overall assembling operation can be further enhanced."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a circuit host card with modules connected thereto and, more particularly, to a circuit host card having three electronic circuit board modules or mezzanine cards connected thereto.\n2. Background Description\nPrinted circuit boards (PCB) or printed wiring boards (PWB) have been known for many years and may be formed by many techniques (e.g. screening, plating, etching, etc.). Printed wiring boards provide for a compact, structurally robust and easily manufactured electronic circuit structure. In use, discrete electronic components and integrated circuits are affixed to the boards via a host of commonly known techniques. Input/Output (I/O) pins are connected to the printed circuit boards (PCB) or printed wiring boards (PWB) in order to make connections to other components of an integrated system.\nAs electronic systems, such as communication systems and data processors, have become more complex, the use of multiple printed circuit boards connected to each other has been implemented. For ease of maintenance and repair, these complex systems have been fabricated in a modular form, often of standardized dimensions. These standardized dimensions, though, place limits on the size and shape of the module and hence the packaging density or packaging area of the system.\nAt the present time, a prevalent standard for modular circuit packaging is commonly known as the xe2x80x9cVersa Module Europaxe2x80x9d (VME) which incorporates a number of design standards including those known in the art as IEEE standard 1101.1, IEEE standard 1101.2, VITA20, ANSI/VITA1, VITA1.1 (VME64X) and IEEE standard P1386. While these standards allow some degree of freedom in component layout within the module, all of the standards involved in the collective VME standard are directed to modules including only a single printed wiring board having components mounted on only a single side thereof. Thus, in order to augment the area of the single VME standard board (also known as a host or main board), one or two circuit boards may be mounted to the host board via interconnections. These boards are known in the art as mezzanine boards or cards. Although the mezzanine boards increase the packaging density, the mezzanine board solution is limited to less than doubling the available area of the host board due to the configuration of the host board and placement of the connectors in accordance with the VME standard.\nAt present, there is only very limited ways known to further augment the printed wiring board area. This is due to several different reasons; namely, for example:\n(i) The size of the host board;\n(ii) The size of the mezzanine boards;\n(iii) The placement and configuration of the I/O pins;\n(iv) The placement and configuration of the connector assemblies between the host board and the mezzanine boards; or\n(v) The cooling requirements of the entire system.\nBeing more specific, the host board is limited to certain dimensions, i.e., approximately 233.5 mm by 160 mm (equivalent to approximately 9xe2x80x3xc3x976.25xe2x80x3), which limits the mounting area of the connectors and mezzanine boards. The mezzanine board is also limited in size to approximately 74 mmxc3x97145 mm (approximately 3xe2x80x3xc3x975.5xe2x80x3). The I/O pins are placed at the edge of the board and are typically designated P1, P2, P0, where P0, P1 and P2 are multiple row multiple-pin configurations. This pin configuration also limits the placement of additional components and mezzanine boards. The mezzanine connector assemblies, on the other hand, may range from two to four connectors, but a first set of connectors must be placed within a range of approximately 40 mm+/xe2x88x9225.5 mnm from a side edge of the host board (leaving a space on the other opposite edge side of the board). This latter requirement provides a very exacting standard for the placement of the mezzanine boards thus, again, limiting the placement of additional components and mezzanine boards. Lastly, cooling ribs run along the length of the host board for dissipating heat generated from the mezzanine boards. This further limits the space and placement of the mezzanine boards. Thus, taking all of this into account, further augmenting the printed wiring board area has been difficult to achieve by card designers.\nBut, to overcome many shortcomings associated with the industry standard, the card design industry has devised several different solutions from adding additional host boards to removing certain functions from the host board or mezzanine boards. The former solution, though, requires additional costs, manufacturing and installation; whereas, the latter solution simply is not a very viable option. Another approach has been to add additional circuit cards to the backplane of the host board. This option adds complicated design considerations such as, for example, connection problems, the placement of the I/O pins and the like. All of these design considerations also add needlessly to the manufacturing and assembly costs.\nThe present invention is directed to overcoming one or more of the problems as set forth above.\nIn a first aspect of the invention, a circuit card assembly includes a host card having connector assemblies and a conduction-cooling path mounted thereon. At least one base mezzanine card having electronic circuitry defining a component field is provided. The at least one base mezzanine card is mounted to the conduction-cooling path of the host card and is further connected to the connector assemblies. A supplementary mezzanine card having electronic circuitry mounted thereon is mounted to the host card and over at least a portion of the component field. The mounting configuration provides a cooling path between the supplementary mezzanine card and the host card and also prevents mechanical interference between the supplementary mezzanine card and the component field. The supplementary mezzanine card is connected to another of the connector assemblies of the host card outside of the connector and PWB component field area of the at least one base mezzanine card.\nIn embodiments, the at least one base mezzanine card is a first base mezzanine card and a second base mezzanine card mounted in substantially a same mounting plane with respect to one another. The supplementary mezzanine card is a third mezzanine card mounted outside the connector and PWB component field area of the first and second base mezzanine cards. The third mezzanine card may be mounted over one of:\n(i) a portion of the component field of the first base mezzanine card,\n(ii) a portion of the component field of the first base mezzanine card and the component field of the second base mezzanine card, and\n(iii) the component field of the first base mezzanine card and the component field of the second base mezzanine card.\nThe circuit card assembly, in further embodiments, includes a center cooling rib running along a length of the host card and side edge cooling ribs running along opposing edges of the host card. The at least one base mezzanine card is mounted to the center cooling rib and the supplementary mezzanine card is mounted to the side edge cooling ribs. In still further embodiments, the supplementary mezzanine card includes a cooling rib corresponding to a location of and mounted to the side edge cooling ribs and a faceplate cooling rib of the host card. A profile of the side edge cooling ribs and the cooling rib prevents mechanical interference between the another mezzanine card and the component field and further provides the cooling path. In yet another embodiment, the supplementary mezzanine card includes opposing cooling ribs corresponding to side edges of the host card for mounting thereon, and has a profile that prevents mechanical interference and provides the cooling path between the another mezzanine card and the host card.\nIn another aspect of the present invention, the host card has connector assemblies and a conduction-cooling path mounted thereon. A first base mezzanine card has electronic circuitry defining a component field. The first base mezzanine card is mounted to the conduction-cooling path and is further connected to a first set of the host mezzanine connector assemblies. A second base mezzanine card has electronic circuitry defining a component field and is mounted to the conduction-cooling path and is connected to a second set of the connector assemblies. A third mezzanine card is mounted to the host card over at least a portion of the component field of the first and/or second base mezzanine cards such that mechanical interference between the third mezzanine card and the component field of the first and second base mezzanine cards is prevented. The third mezzanine card is further connected to a third set of the connector assemblies of the host card outside of the connecting area and component field of the first and second mezzanine cards. At least the conduction-cooling path provides a cooling path between the first and second mezzanine cards and the host card.\nIn embodiments, the conduction-cooling path of the host card also provides a cooling path between the third mezzanine card and the host card. Alternatively, a conduction-cooling path on either the third mezzanine card or both the third mezzanine card and the host card provides the cooling path.\nIn yet another embodiment, a circuit card assembly includes first, second and third mezzanine cards mounted and connected to a host card. The third mezzanine card is connected outside a connecting area of the first and second mezzanine cards. A cooling path of the host card is used to mount the first, second and third mezzanine cards. The third mezzanine card also includes a conduction-cooling path that is mounted to the conduction-cooling path of the host card. The combination of the conduction-cooling paths provide a cooling path between the third mezzanine card and the host card. In embodiments, the profile of the conduction-cooling paths also prevents mechanical interference between the third mezzanine card and the first and second mezzanine cards."}
{"text": "This invention relates to systems and methods for processing digital images. More specifically it relates to image segmentation wherein an object of interest in an image is separated from the background of that object in the image.\nSegmentation is a frequently encountered problem in image processing. It is particularly important in medical imaging where it is desired to extract an object, such as an organ or tumor, from an image for further processing.\nThere are many different segmentation processes available. The effectiveness of all of the segmentation processes vary in accordance with the quality of the image data, the application as well as a number of other factors.\nSegmentation in medical image data is particularly difficult. The massive amounts of data make segmentation difficult in this area. Further, the quality of the image data obtained from medical imaging devices is not always optimal. Additionally, the quality of the segmentation is generally very important since the segmentation may be used to guide a diagnosis or a treatment plan.\nAccordingly, new and improved methods and systems are required that will create segmentation of images and image volumes."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor integrated circuit and method of testing the same.\n2. Description of Related Art\nIn recent years, as an LSI (Large Scale Integration) has gotten larger in scale, the number of PLLs (Phase-Locked Loops) embedded in an LSI has been increasing. Consequently, the testing time of PLLs embedded in a LSI has been problematically getting longer.\nA conventional method of testing PLLs embedded in a LSI is explained hereinafter with reference to FIG. 8. In FIG. 8 shows an exemplary case where a LSI 1000 has two PLLs, i.e., first PLL 1001 and second PLL 1002. A test device 2000 for testing a PLL includes a signal generator 2001 and a signal measuring device 2002. The test device 2000 also includes switches 2003 and 2004. The switch 2003 changes the connection from the signal generator 2001 between the first PLL 1001 and second PLL 1002. Likewise, the switch 2004 changes the connection to the signal measuring device 2002 between the first PLL 1001 and second PLL 1002.\nThen, to test the PLLs embedded in the LSI 1000, firstly, the first PLL 1001 is connected to the signal generator 2001 and signal measuring device 2002 by the first and second switches 2003 and 2004. In this state, a clock having a frequency ft outputted from the signal generator 2001 is inputted to the first PLL 1001, and the frequency is frequency-multiplied by N at the first PLL 1001. Then, the clock having a frequency N×ft outputted from the first PLL 1001 is measured in the signal measuring device 2002, and the test of the first PLL 1001 has been completed. Next, by changing the switches 2003 and 2004, the second PLL 1002 is connected to the signal generator 2001 and signal measuring device 2002. In this state, a clock having a frequency ft outputted from the signal generator 2001 is inputted to the second PLL 1002, and the frequency is frequency-multiplied by M at the second PLL 1002. Then, the clock having a frequency M×ft outputted from the second PLL 1002 is measured in the signal measuring device 2002, and the test of the second PLL 1002 has been completed.\nFIG. 9 shows a more concrete example of the circuit shown in FIG. 8. As shown in FIG. 9, selectors 1004 and 1005 are connected to the first and second PLLs 1001 and 1002 respectively. In a normal operation mode, the selectors 1004 and 1005 select a clock generated by an OSC (oscillator) 1003, and this clock is inputted to the first and second PLLs 1001 and 1002. On the other hand, in a test mode, the selectors 1004 and 1005 select a clock generated by the signal generator 2001, and this clock is inputted to the first and second PLLs 1001 and 1002. Furthermore, first and second logic circuits 1006 and 1007 are connected to the first and second PLLs 1001 and 1002. The first and second logic circuits 1006 and 1007 become active with the clock outputted from the first and second PLLs 1001 and 1002 in the normal mode, Furthermore, the test device 2000 outputs a control signal for controlling the switching of the selector 1004 and 1005.\nThen, to test the PLLs embedded in the LSI 1000, firstly, the first PLL 1001 is connected to the signal generator 2001 by a switch 2005, and the first PLL 1001 is also connected to the signal measuring device 2002 by a switch 2006. At the same time, the test device 2000 inputs the control signal to the selector 1004 such that a clock generated by the signal generator 2001 is inputted to the first PLL 1001. In this state, a clock having a frequency ft outputted from the signal generator 2001 is inputted to the first PLL 1001, and the frequency is frequency-multiplied by N at the first PLL 1001. Then, the clock having a frequency N×ft outputted from the first PLL 1001 is measured in the signal measuring device 2002, and the test of the first PLL 1001 has been completed. Next, the second PLL 1002 is connected to the signal generator 2001 by changing the switch 2005, and the second PLL 1002 is also connected to the signal measuring device 2002 by changing the switch 2006. At the same time, the test device 2000 inputs the control signal to the selector 1005 such that a clock generated by the signal generator 2001 is inputted to the second PLL 1002. In this state, a clock having a frequency ft outputted from the signal generator 2001 is inputted to the second PLL 1002, and the frequency is frequency-multiplied by M at the second PLL 1002. Then, the clock having a frequency M×ft outputted from the second PLL 1002 is measured in the signal measuring device 2002, and the test of the second PLL 1002 has been completed.\nHowever, in the method shown in FIGS. 8 and 9, in the case where a plurality of PLLs are embedded in a LSI, a test has to be performed the same number of times as the number of the PLLs embedded in the LSI to complete the PLL test. Therefore, the testing time of the PLLs gets problematically longer. Furthermore, since both a signal generator and a signal measuring device are expensive, it will be unrealistic to use several test devices simultaneously. In addition, the number of signal generators and signal measuring devices embedded in a single test device is two channels or three at the maximum, and even so it cannot test more than two PLLs simultaneously.\nMeanwhile, it is known to test two PLLs simultaneously for a LSI having two PLLs by inputting a clock delayed by a delay circuit to one of the PLLs, comparing the outputted clocks from the two PLLs by a comparator, and detecting the failure of the PLLs from the phase difference between the outputted clocks from the two PLLs. (For example, Japanese Unexamined Patent Application Publication No. 2005-277472 (Ogawa))\nHowever, the technique described in Ogawa can only compare the phase difference between two input clocks, and cannot cope with a LSI having more than two PLLs."}
{"text": "The structure and operation of hard disk drives is generally known. Hard disk drives include, generally, a case, a hard disk having magnetically alterable properties, and a read/write mechanism including Read/Write (RW) heads operable to write data to the hard disk by locally alerting the magnetic properties of the hard disk and to read data from the hard disk by reading local magnetic properties of the hard disk. The hard disk may include multiple platters, each platter being a planar disk.\nAll information stored on the hard disk is recorded in tracks, which are concentric circles organized on the surface of the platters. FIG. 1 depicts a pattern of radially-spaced concentric data tracks 12 within a disk 10. Data stored on the disks may be accessed by moving RW heads radially as driven by a head actuator to the radial location of the track containing the data. To efficiently and quickly access this data, fine control of RW head positioning is required. The track-based organization of data on the hard disk(s) allows for easy access to any part of the disk, which is why hard disk drives are called “random access” storage devices.\nSince each track typically holds many thousands of bytes of data, the tracks are further divided into smaller units called sectors. This reduces the amount of space wasted by small files. Each sector holds 512 bytes of user data, plus as many as a few dozen additional bytes used for internal drive control and for error detection and correction.\nTypically, these tracks and sectors are created during the low level formatting of the disk. This low level formatting process creates the physical structures (tracks, sectors, control information) on the disk. Normally, this step begins with the hard disk platters containing no information. Newer disks use many complex internal structures, including zoned bit recording to put more sectors on the outer tracks than the inner ones, and embedded servo data to control the head actuator. Newer disks also transparently map out bad sectors. Due to this complexity, all modern hard disks are low-level formatted at the factory for the life of the drive.\nThe ability to store and access increased amounts of data depends on the ability to accurately position the RW head relative to the data tracks. Positioning of the RW head relative to the physical structures is typically based on amplitude information of two (2) or four (4) bursts within a servo pattern. This amplitude information is subject to noise, and environmental changes. A four burst amplitude servo pattern is illustrated in FIG. 2. For example, when the RW right head is located over center of the N−1 track of FIG. 2 a maximum “A” signal amplitude, minimum “B” signal amplitude, and “C” and “D” signal amplitudes that are approximately equal are sensed. The amplitude of bursts “A”, and “B”, must allow the disk controller to determine whether the head is over an even or odd track while the “C” and “D” signal amplitudes allow the controller to determine the distance how well centered the head is relative to the track. The amplitudes of signals “A”, “B”, “C”, and “D”, associated with varying track position is shown in FIG. 3A.\nIn FIG. 3B a servo signal is shown on an individual track. As can be seen the servo signal contains a preamble sync mark, and gray code information prior to the a four bursts, this case the amplitudes of the “A” and “B” signal bursts are approximately equal. The “C” burst is at a maximum while the “D” burst is at a minimum. Returning to FIG. 2 this would correspond to the RW head being located midway between track N and N−1 where the “A” and “B” amplitudes are approximately equal. The “C” amplitude signal is at a maximum while the “D” amplitude signal is at a minimum. The data portion of the servo signal follows the four bursts.\nThis four burst pattern does not allow all the disk space to be effectively used. This servo pattern requires unused media leaves less user data space on the media for storage. To realize additional storage availability, require higher data density is required which may result in a poor quality and production yield.\nFurther limitations and disadvantages of conventional and traditional RW head positioning processes and related functionality will become apparent to one of ordinary skill in the art through comparison with the present invention described herein."}
{"text": "Long-read sequencing and long-range mapping technologies may be used to generate more accurate genome assemblies. In short, there are many genomic regions, and many kinds of structural variation, that cannot be correctly assembled from short-read paired-end sequencing data, including typical sequencing data produced by the industry leaders (such as Illumina and Thermo Fisher Life Technologies). However, many of these regions and variations can be correctly assembled using newer technologies that generate primary read lengths (or genomic maps) measured in the 10's to 100's of kilobases (kb), such as those technologies developed by Pacific Biosciences, Oxford Nanopore, 10X Genomics, Genomic Vision, Roche/Genia, and Bionano Genomics.\nAlthough the long-read technologies are still in an early stage of development, they may be used to complement short-read sequencing capabilities. For example, DNA extraction and library preparation technologies can be used to produce long and high-quality DNA libraries. However, a gap in the marketplace exists for commercial technologies that can produce long DNA fragments. Most commercial genomic DNA extraction kits yield maximum DNA size of 20-50 kb. Since some current systems are capable of primary reads greater than 50 kb, however, existing DNA extraction kits limit the capabilities of these systems. Additionally, long DNA fragments may be used for optical mapping and synthetic long-read systems (e.g., 10X Genomics, Bionano Genomics), which require libraries generated from genomic DNA samples that are 100's of kb in length. Therefore, there is a need for automated and reproducible methods for producing extremely long (e.g., 100's to 1000's of kb in length), high quality genomic DNA samples. Despite a demonstrated need and efforts to achieve these samples, a solution has not yet been developed.\nIn addition, library preparation is a multi-step process that is divided into two major divisions, 1) DNA extraction from biological samples (biological fluids, particulates, cells, and tissue), and 2) library construction. As DNA sequencing becomes more useful to clinical studies, diagnostics, and therapy management, integrated workflows may streamline and automate the overall process. Again, although there is a demonstrated need, increased integration of the overall sample to library workflow has not yet been achieved.\nFurthermore, next-generation sequencing (NGS) for detecting and diagnosis of infectious disease has great potential in molecular diagnostics. For example, clinical workflows using clinical samples (frequently whole blood or white blood cells (WBC's) from whole blood) may be subjected to NGS sequencing, followed by in-silico subtraction of all sequences that can be mapped to the human reference sequence. The remaining “non-human” sequences are examined for identity with any known pathogen sequences. Although this process may be successful, it requires very efficient library construction methods and deep sequencing runs, both of which can be expensive. This expense has kept the in-silico subtraction method from widespread use."}
{"text": "Cancer is one of the leading causes of premature death worldwide. Furthermore, colorectal cancer is the second leading cause of death from cancer in the western world. Colorectal cancer is curable if diagnosed early. One of the major treatment protocols for colorectal cancer is surgery. With no additional therapy, surgery typically cures Duke's A colon cancer that invades the submucosa and Duke's B1 disease that invades the muscularis propria (Cohen et al., In: DeVita et al. (eds.), Cancer Principles and Practice of Oncology, Ed. 5, pp. 1144-1197. New York: Lippincott-Raven.).\nThe primary prognostic approach to identify differences among patients in early stages of the disease is the Tumor-Node-Metastasis (TNM) system (Graziano et al. (2003) Ann Oncol. 14, 1026-1038). However, the survival outcome varies among patients with similar pathological disease stages. There have been increasing demands to specify the molecular markers for more aggressive colorectal cancer in order to identify the disease well in advance and prescribe appropriate patient therapy. In response, new therapeutic strategies are needed to combat colorectal cancer.\nIt is known that the activity of N-myristoyltransferase (NMT) is highly elevated in colorectal cancer (Magnuson et al. 1995, J. Natl. Cancer Inst. 87, 1630-1635). NMT is an enzyme that catalyzes the myristoylation of proteins involved in diverse biological functions including oncogenesis (Rajala et al. (2000) Mol. Cell. Biochem. 204, 135-155; Selvakumar et al., (2002) Int. J. Mol. Med. 10, 493-500).\nN-myristoyltransferase activity is increased in rat colonic tumors compared to those of normal appearing colonic mucosa (Magnuson et al., (1995) J. Natl. Cancer Inst. 87, 1630-1635.). Interestingly, NMT activity in the normal appearing mucosa of colon cancer bearing rats was similar to the activity found in colonic mucosa of control rats, suggesting that elevated NMT activity is restricted to the tumor site. Higher NMT activity is observed in rat colonic tumors when compared with the corresponding normal-appearing or normal mucosa. Furthermore, a several fold increase in NMT activity as compared to adjacent normal-appearing mucosa, was observed in polyps and stage B1 tumors located in the descending colon close to the rectum.\nIn addition to the rat model, elevated NMT activity and expression was also observed in human adenocarcinoma compared to that of normal appearing mucosa (Magnuson et al., (1995) J. Natl. Cancer Inst. 87, 1630-1635; Raju et al., (1997) Exp. Cell Res. 235, 145-154.). NMT activity in normal appearing mucosa is similar to that of Crohn's disease and the volvulus, indicating that elevated NMT activity is specific to cancer and is not a non-specific response to inflammatory conditions or noncancerous lesions. Rajala et al. (Cancer 2000; 88(9):1992-1999) showed increased expression on NMT in gall bladder tissue of patients with gall bladder cancer.\nPrevious work, however, has not established NMT as a marker for cancer, given the limitation that the elevated expression and activity of NMT in the tumor region is not accessible for prognostic/diagnostic purposes until a colonoscopy is performed and a tissue sample is available for protein analysis.\nThere is a need for non-invasive or minimally invasive technologies that can be used to detect cancer."}
{"text": "The present disclosure relates to a solid-wall screw centrifuge including a discharge for discharging a clarified liquid from a drum of the centrifuge. The drum includes a passage in a drum cover associated with a weir plate. The passage, drum cover, and weir plate form an overflow weir on a radial inner edge. A deflecting device is arranged on an outside on the drum cover. The deflecting device is arranged such that an exiting fluid flow from the centrifuge is one of deflected in the circumferential direction and subjected to a twist in the circumferential direction.\nSolid-wall screw centrifuges are known.\nDE 102 03 652 B4 discloses a solid-wall screw centrifuge comprising a device for discharging clarified fluid from a drum having a drum cover with a passage. The passage associated with a throttle device, for example, an orifice gauge, the distance of which to the passage is variable. The passage is further associated with nozzles for the discharge of the clarified fluid, which nozzles are aligned obliquely in relation to the radial direction through the drum axis for saving energy.\nNozzles on solid-wall screw centrifuges and their effect of energy saving in a respective alignment obliquely in relation to the drum axis are further known from DE 39 04 151 A1.\nThe publication of “Patent Abstracts of Japan”, No. 11179236A, discloses that guide plates can be associated with passage openings in the drum cover which are respectively partly sealed by a weir plate forming an overflow weir. The guide plates provide fluid exiting the drum with a twist, wherein the occurring repulsion effect is to be used for energy saving. The guide plates are attached to the outside of the overflow weir on the weir plates, for example. They are arranged as flat plates, for example, which are aligned parallel to the rotational axis of the drum. The plane in which the flat plates are disposed does not intersect the drum axis and encloses an angle of approximately 45° with the radially extending straight line which, respectively, extends through the rotational axis of the drum and the respective passage opening.\nThe present disclosure thus relates to embodiments that provide a solid-wall screw centrifuge which is optimized with respect to energy saving.\nThe present disclosure further relates to a solid-wall screw centrifuge that includes a discharge for discharging a clarified liquid from a drum of the centrifuge. The drum includes a passage in a drum cover associated with a weir plate, which passage, drum cover, and weir plate form an overflow weir on a radial inner edge. A deflecting device is arranged on an outside on the drum cover, the deflecting device being arranged such that an exiting fluid flow from the centrifuge is one of deflected in the circumferential direction and subjected to a twist in the circumferential direction. The deflecting device is arranged so as to maintain the position of the overflow weir such that the exiting fluid flow is deflected in the circumferential direction at an overflow edge of the weir flute and directly in an axial extension of a region of a passage opening that is not covered by the weir plate\nThe deflecting device is, accordingly, arranged in such a way, for example, on the weir plate, that, by maintaining the overflow weir, an exiting fluid flow is subjected to a twist in the circumferential direction at the overflow edge in the axial extension of the passage opening.\nThis can be achieved, according to an embodiment of the present disclosure, in such a way, for example, that the deflecting wall is disposed partly directly in axial extension of the region of the passage opening not covered by the respective weir plate axially outside before the passage opening, and axially partly or fully covers the opening.\nIt is also known, for example, from U.S. 2004/0072668 A1 or WO 2008/138345 A1,to arrange a housing instead of an overflow weir on the weir plate of the passage opening, with the overflow weir or a nozzle being then arranged in an oblique side of the housing on the passage opening. These constructions are disadvantageous, however, in that co-rotating fluid chambers are formed on the outside of the drum cover, which fluid chambers form a part of the rotating system, so that the clarification or separation process can still continue in the fluid chambers. This separation process may lead to undesirable deposits in the chambers which lead to cleaning problems.\nIt is, therefore, advantageous that an embodiment according to the present disclosure leaves the overflow weir directly on the weir plates on the drum cover, and only the fluid which has already left the rotating system via the overflow weir will be deflected in the circumferential direction. After the exit from the drum, no clarification effect can occur any more. Instead, the entire exiting fluid jet will rather be deflected by the deflecting device in the circumferential direction.\nAdvantageous embodiments of the present disclosure are discussed herein and in the appended claims.\nOther aspects of the present disclosure will become apparent from the following descriptions when considered in conjunction with the accompanying drawings."}
{"text": "The present invention pertains to a transport or conveying assembly for a harvest recovery machine, the assembly including a rotor, at least one stripper and a crop transport channel that includes at least one guide wall.\nDE-C1 42 19 719 discloses a baling press for large bales which includes a transport channel through which the harvested crop material is conveyed from a pick-up to a compression channel by a stuffing rotor. The transport channel includes a guide wall that is formed of webs that extend in the material flow direction and guide plates situated between these webs. The guide plates are able to move against the force of springs in order to prevent peak loads.\nDE-C2 195 46 263 discloses a transport device for a round baler, in which the transport channel includes a lower guide wall that can be pivoted transverse to the material flow direction about a downstream bearing. A movement of the guide wall and consequently a widening of the inlet gap occurs when a certain pressure of the harvest material is reached.\nAccording to DE-C1 198 21 591, the bottom of a transport section between a pick-up and the compression chamber of a baler is held in a pivot bearing on its upstream end and adjustable at a downstream end region, namely transverse to the material flow. The adjustment is realized with the aid of a motor and serves for removing obstructions from the transport channel.\nAccording to EP-A1 0 074 533 a guide wall of a transport device of a self-loading forage box can be adjusted with the aid of a lever in order to remove obstructions.\nEP-A1 0 339 733 discloses a pivoted wall that makes it possible to open the transport channel in order to make it possible to remove obstructions from the supply channel of a large agricultural baler.\nAccording to DE-A1 198 41 598, a round baler is provided with a cutting mechanism in combination with a rotor. Knives that protrude through a wall and extend into a transport channel are provided opposite the rotor. The wall can be adjusted transverse to the material flow upstream of the rotor in order to change the channel cross section. On the downstream end, the wall can be pivoted downward in order to remove foreign matter from the channel. The harvest material is removed from the rotor at the harvest material inlet by a stationary stripper. The adjustment of the wall takes place in dependence on a torque that is measured on the rotor.\nThe problem on which the invention is based can be seen in the fact that known measures aim to remove or eliminate rather than prevent obstructions or peak loads or require a significant adjustment of the guide wall.\nAccording to the present invention there is provided an improved arrangement of adjustable components for transporting harvest material from a pickup into a baling chamber.\nAn object of the invention is to provide components arranged for defining a channel through which harvest material is transported and for effecting the transport of such material from a pick-up to the inlet of an arrangement for further processing, with at least some of the components being adjustable to vary the cross section of the channel in accordance with the amount of crop so as to maintain a constant density.\nA more specific object of the invention is to provide transport components, as set forth in the previous object, wherein one of the components is a guide wall which is adjustable transversely to the flow path of the material as it progresses from the pick-up to the inlet.\nAnother object of the invention is to provide transport components, as set forth in one or more of the previous objects, wherein a tined rotor is provided for acting in conjunction with the guide wall to form a portion of the channel, with the tined rotor being adjustable alone or together with the guide wall in a direction transverse to the flow path of the material as it progresses from the pick-up to the inlet.\nStill another object of the invention is to provide transport components, as set forth in one or more of the previous objects, wherein a stripper assembly is associated with the tined rotor and acting in conjunction with the guide wall to form a portion of the channel, with the stripper assembly being adjustable alone or together with the guide wall and/or the rotor in a direction transverse to the flow path of the material as it progresses from the pick-up to the inlet.\nThese and other objects will become apparent from a reading of the ensuing description together with the appended drawings."}
{"text": "This invention relates generally to field effect transistors (FETs), and more particularly the invention relates to use of a stressed capping layer to enhance the performance of complementary metal-oxide-semiconductor (CMOS) integrated-circuit (IC) devices.\nCMOS IC devices comprise n-channel (NMOS) and p-channel (PMOS) transistors which toggle on-off in a complementary fashion in response to input voltage signals. The transistors can have a planar structure, in which the surfaces of the source, channel, and drain regions are substantially located in a single plane, as in a conventional bulk-silicon (bulk-Si) MOSFET or silicon-on-insulator (SOI) MOSFET, or the transistors can have a vertical structure such as in the double gate FinFET structure as described in U.S. Pat. No. 6,413,802.\nTransistor performance depends on mobile carrier (hole or electron) mobility. Strained-silicon technologies have been widely investigated to enhance the performance of CMOS devices. In particular, strain induced by the use of a stressed SiNx capping layer is advantageous because of its process simplicity and its extendibility from bulk-Si to SOI MOSFETs. However, it is difficult to enhance the performance of both n-channel and p-channel devices simultaneously, because tensile stress is required in the n-channel MOSFET channel, whereas uniaxial compressive stress is required in the p-channel MOSFET channel, to enhance carrier mobility.\nThe present invention is directed to the fabrication of a CMOS device in which the performance of both n-channel and p-channel transistors can be enhanced with a common capping layer and to the resulting structure."}
{"text": "1. Technical Field\nThe disclosure relates to an electronic apparatus and particularly relates to an electronic apparatus which uses a pivoting-locking module and a supporting frame to achieve the locking and pivoting between two bodies.\n2. Background\nAs technology advances and develops, various kinds of electronic devices, such as mobile phones, tablet computers, and electronic books, etc., have been extensively used. These electronic devices make people's life more convenient and have become an indispensable part to people's life. These electronic devices are usually made small and light for the users to use them with hands and to carry them more easily.\nTake the tablet computer as an example, when inputting text, the user may use a virtual keyboard displayed on the touch screen of the tablet computer to input text via touch control. However, touch control is inconvenient and inefficient for quick input of text. In that case, an extended physical keyboard may be applied to solve the problem."}
{"text": "Automatic processor frequency governors are software components within an operating system (“OS”) that optimize a platform's power/performance by automatically scaling the processor (“CPU”) frequency based on CPU load. Typically, the governor analyzes processor utilization at regular intervals, such as 250 ms, dropping the frequency in response to high-idle periods, and raising it in response to high-utilization periods. CPU utilization is typically calculated by detecting if the CPU is halted when a periodic timer interrupt on the platform is fired. Thus, for example, if the timer interrupt fires 1000 times per second and in the last 1/10 of a second, the CPU was halted 40 of 100 times when the timer interrupt fired, the OS may calculate a CPU “load” of 60%."}
{"text": "The present invention relates generally to recombinant DNA technology; and more specifically, to the development of recombinant vectors useful for directing the expression of one or more heterologous gene products.\nAlphaviruses comprise a set of genetically, structurally, and serologically related arthropod-borne viruses of the Togaviridae family. These viruses are distributed worldwide, and persist in nature through a mosquito to vertebrate cycle. Birds, rodents, horses, primates, and humans are among the defined alphavirus vertebrate reservoir/hosts.\nTwenty-six known viruses and virus subtypes have been classified within the alphavirus genus utilizing the hemagglutination inhibition (HI) assay. This assay segregates the 26 alphaviruses into three major complexes: the Venezuelan equine encephalitis (VEE) complex, the Semliki Forest (SF) complex, and the western equine encephalitis (WEE) complex. In addition, four other viruses, eastern equine encephalitis (EEE), Barmah Forest, Middelburg, and Ndumu, receive individual classification based on the HI serological assay.\nMembers of the alphavirus genus also are classified based on their relative clinical features in humans: alphaviruses associated primarily with encephalitis, and alphaviruses associated primarily with fever, rash, and polyarthritis. Included in the former group are the VEE and WEE complexes, and EEE. In general, infection with this group can result in permanent sequelae, including behavior changes and learning disabilities, or death. In the latter group is the SF complex, comprised of the individual alphaviruses Semliki Forest, Sindbis, Ross River, Chikungunya, O\"\"nyong-nyong, and Mayaro. With respect to this group, although serious epidemics have been reported, infection is in general self-limiting, without permanent sequelae.\nSindbis virus is the prototype member of the Alphavirus genus of the Togaviridae family. Its replication strategy after infection of cells (see FIG. 1) has been well characterized in chicken embryo fibroblasts (CEF) and baby hamster kidney (BHK) cells, where Sindbis virus grows rapidly and to high titer, and serves as a model for other alphaviruses. Briefly, the genome from Sindbis virus (like other alphaviruses) is an approximately 12 kb single-stranded positive-sense RNA molecule which is capped and polyadenylated, and contained within a virus-encoded capsid protein shell. The nucleocapsid is further surrounded by a host-derived lipid envelope into which two viral-specific glycoproteins, E1 and E2, are inserted and anchored to the nucleocapsid. Certain alphaviruses (e.g., SF) also maintain an additional protein, E3, which is a cleavage product of the E2 precursor protein, PE2. After virus particle absorption to target cells, penetration, and uncoating of the nucleocapsid to release viral genomic RNA into the cytoplasm, the replicative process is initiated by translation of the nonstructural proteins (nsPs) from the 5xe2x80x2 two-thirds of the viral genome. The four nsPs (nsP1-nsP4) are translated directly from the genomic RNA template as one of two polyproteins (nsP123 or nsP1234), and processed post-translationally into monomeric units by an active protease in the C-terminal domain nsP2. A leaky opal (UGA) codon present between nsP3 and nsP4 of most alphaviruses accounts for a 10 to 20% abundance of the nsP1234 polyprotein, as compared to the nsP123 polyprotein. Both of the nonstructural polyproteins and their derived monomeric units may participate in the RNA replicative process, which involves binding to the conserved nucleotide sequence elements (CSEs) present at the 5xe2x80x2 and 3xe2x80x2 ends, and a junction region subgenomic promoter located internally in the genome (discussed further below).\nThe positive strand genomic RNA serves as template for the nsP-catalyzed synthesis of a full-length complementary negative strand. Synthesis of the complementary negative strand is catalyzed after binding of the nsP complex to the 3xe2x80x2 terminal CSE of the positive strand genomic RNA. The negative strand, in turn, serves as template for the synthesis of additional positive strand genomic RNA and an abundantly expressed 26S subgenomic RNA, initiated internally at the junction region promoter. Synthesis of additional positive strand genomic RNA occurs after binding of the nsP complex to the 3xe2x80x2 terminal CSE of the complementary negative strand genomic RNA template. Synthesis of the subgenomic mRNA from the negative strand genomic RNA template, is initiated from the junction region promoter. Thus, the 5xe2x80x2 end and junction region CSEs of the positive strand genomic RNA are functional only after they are transcribed into the negative strand genomic RNA complement (i.e., the 5xe2x80x2 end CSE is functional when it is the 3xe2x80x2 end of the genomic negative stranded complement). The structural proteins (sPs) are translated from the subgenomic 26S RNA, which represents the 3xe2x80x2 one-third of the genome, and like the nsPs, are processed post-translationally into the individual proteins.\nSeveral groups have suggested utilizing certain members of the alphavirus genus as an expression vector, including, for example, Sindbis virus (Xiong et al., Science 243:1188-1191,1989; Hahn et al., Proc. Natl. Acad. Sci. USA 89:2679-2683, 1992; Dubensky et al., J. Virol. 70:508-519,1996), Semliki Forest virus (Liljestrom, Bio/Technology 9:1356-1361,1991), and Venezuelan Equine Encephalitis virus (Davis et al., J. Cell. Biochem. Suppl. 19A:10,1995). In addition, one group has suggested using alphavirus-derived vectors for the delivery of therapeutic genes in vivo. One difficulty, however, with the above-referenced vectors is that inhibition of host cell-directed macromolecular synthesis (i.e., protein or RNA synthesis) begins within a few hours after infection and cytopathic effects (CPE) occur within 12 to 16 hours post infection (hpi). Inhibition and shutoff of host cell protein synthesis begins within 2 hpi in BHK cells infected with recombinant viral particles, in the presence or absence of structural protein expression, suggesting that the early events after virus infection (e.g., synthesis of nsPs and minus strand RNA) may directly influence the inhibition of host cell protein synthesis and subsequent development of CPE and cell death.\nSIN-1 is a variant strain derived from wild-type Sindbis, and was isolated from a culture of BHK cells persistently infected with Sindbis virus over a period of one month (Weiss et al. J. Virol. 33: 463-474. 1980). A pure SIN-1 virus stock obtained by expansion from a single plaque does not kill the BHK cells which it infects. Importantly, virus yields ( greater than 103 PFU/cell) are the same in BHK cells infected with wild-type Sindbis virus or the variant SIN-1 virus. Thus, the principle phenotype of SIN-1 in infected BHK cells is characterized by production of wild-type levels of infectious virus in the absence of virus-induced cell death.\nThe present invention provides recombinant vectors with selected desirable phenotypes for use in a variety of applications, including for example, gene therapy and recombinant protein production, and further provides other related advantages.\nBriefly stated, the present invention provides RNA vector replicons, alphavirus vector constructs, eukaryotic layered vector initiation systems and recombinant alphavirus particles which exhibit reduced, delayed, or no inhibition of cellular macromolecular synthesis (e.g., protein or RNA synthesis), thereby permitting the use of these vectors for protein expression, gene therapy and the like, with reduced, delayed, or no development of CPE or cell death. Such vectors may be constructed from a wide variety of alphaviruses (e.g., Semliki Forest virus, Ross River virus, Venezuelan equine encephalitis virus or Sindbis virus), and designed to express numerous heterologous sequences (e.g., a sequence corresponding to protein, a sequence corresponding to antisense RNA, a sequence corresponding to non-coding sense RNA, or a sequence corresponding to ribozyme).\nWithin one aspect of the invention, isolated nucleic acid molecules are provided comprising an altered alphavirus nonstructural protein gene which, when operably incorporated into a recombinant alphavirus, increases the time required to reach 50% inhibition of host-cell directed macromolecular synthesis following expression in mammalian cells, as compared to a wild-type alphavirus. As utilized within the context of the present invention, xe2x80x9caltered alphavirus nonstructural protein genexe2x80x9d refers to a gene which, when operably incorporated into an alphavirus RNA vector replicon, recombinant alphavirus particle, or eukaryotic layered vector initiation system, produces the desired phenotype (e.g., reduced, delayed or no inhibition of cellular macromolecular synthesis). Such altered alphavirus nonstructural protein genes will have one or more nucleotide substitutions, deletions, or insertions, which alter the nucleotide sequence from that of the wild-type alphavirus gene. The gene may be derived either artificially (e.g., from directed selection procedures; see Example 2 below), or from naturally occurring viral variants (see Example 1 below). In addition, it should be understood that when the isolated nucleic acid molecules of the present invention are incorporated into an alphavirus RNA vector replicon, recombinant alphavirus particle, or eukaryotic layered vector initiation system as discussed above, that they may, within certain embodiments, substantially increase the time required to reach 50% inhibition of host-cell directed macromolecular synthesis, up to and including substantially no detectable inhibition of host-cell directed macromolecular synthesis (over any period of time). Assays suitable for detecting percent inhibition of host-cell directed macromolecular synthesis include, for example, that described within Example 1.\nWithin other aspects of the invention, isolated nucleic acid molecules are provided comprising an altered alphavirus nonstructural protein gene which, when operably incorporated into a recombinant alphavirus particle, eukaryotic layered vector initiation system, or RNA vector replicon, results in a reduced level (e.g., 2-fold, 5-fold, 10-fold, 50-fold or more than 100-fold) of vector-specific RNA synthesis as compared to the wild-type, and the same or greater level of protein encoded by RNA transcribed from the viral junction region promoter, as compared to a wild-type recombinant alphavirus particle, wild-type eukaryotic layered vector initiation system, or wild-type RNA vector replicon. Representative assays for quantitating RNA levels include [3H] uridine incorporation as described in Example 1, or RNA accumulation as detected by Northern Blot analysis (see Example 4). Representative assays for quantitating protein levels include scanning densitometry (see Example 4) and various enzymatic assays (see Examples 3-5).\nWithin one embodiment of the above, the isolated nucleic acid molecule encodes nonstructural protein 2 (nsP2). Within a further embodiment, the isolated nucleic acid molecule has a mutation in the LXPGG motiff of nsP2.\nWithin another aspect of the invention, expression vectors are provided comprising a promoter operably linked to one of the above-described nucleic acid molecules. Within one embodiment, the expression vector further comprises a polyadenylation sequence or transcription termination sequence 3xe2x80x2 to the nucleic acid molecule.\nWithin yet another aspect of the present invention, alphavirus vector constructs are provided, comprising a 5xe2x80x2 promoter which initiates synthesis of viral RNA in vitro from cDNA, a 5xe2x80x2 sequence which initiates transcription of alphavirus RNA, a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins including an isolated nucleic acid molecule as described above, an alphavirus viral junction region promoter, an alphavirus RNA polymerase recognition sequence and a 3xe2x80x2 polyadenylate tract.\nWithin a related aspect, such constructs further comprise a selected heterologous sequence downstream of and operably linked to a viral junction region. Within a related aspect, alphavirus vector constructs are provided comprising a 5xe2x80x2 promoter which initiates synthesis of viral RNA in vitro from cDNA, a 5 xe2x80x2 sequence which initiates transcription of alphavirus RNA, a nucleic acid molecule which operably encodes all four alphavirus non-structural proteins, an alphavirus viral junction region promoter, an alphavirus RNA polymerase recognition sequence, and a 3xe2x80x2 polyadenylate tract, wherein said in vitro synthesized RNA, upon packaging into an alphavirus particle and introduction of the particle into a mammalian host cell, increases the time required to reach 50% inhibition of host-cell directed macromolecular synthesis following expression in mammalian cells, as compared to a wild-type alphavirus particle.\nWithin a further aspect, alphavirus vector constructs are provided comprising a 5xe2x80x2 promoter which initiates synthesis of viral RNA in vitro from cDNA, a 5xe2x80x2 sequence which initiates transcription of alphavirus RNA, a nucleic acid molecule which operably encodes all four alphavirus non-structural proteins, an alphavirus viral junction region promoter, an alphavirus RNA polymerase recognition sequence, and a 3xe2x80x2 polyadenylate tract, wherein said in vitro synthesized RNA, upon packaging into an alphavirus particle and introduction of the particle into a mammalian host cell, has a reduced level of vector-specific RNA synthesis as compared to wild-type alphavirus particle, and the same or greater level of protein encoded by RNA transcribed from the viral junction region promoter, as compared to a wild-type alphavirus particle.\nWithin yet other aspects of the present invention, RNA vector replicons capable of translation in a eukaryotic system are provided, comprising a 5xe2x80x2 sequence which initiates transcription of alphavirus RNA, a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins, including the isolated nucleic acid molecules discussed above, an alphavirus viral junction region, an alphavirus RNA polymerase recognition sequence and a 3xe2x80x2 polyadenylate tract.\nWithin a related aspect, alphavirus RNA vector replicons capable of translation in a eukaryotic system are provided, comprising a 5xe2x80x2 sequence which initiates transcription of alphavirus RNA, a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins, an alphavirus viral junction region promoter, an alphavirus polymerase recognition sequence and a 3xe2x80x2 polyadenylate tract, wherein said alphavirus RNA, upon packaging into an alphavirus particle and introduction of the particle into a mammalian host cell, increases the time required to reach 50% inhibition of host-cell directed macromolecular synthesis following expression in mammalian cells, as compared to a wild-type alphavirus particle.\nWithin other aspects, alphavirus RNA vector replicons capable of translation in a eukaryotic system are provided comprising a 5xe2x80x2 sequence which initiates transcription of alphavirus RNA, a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins, an alphavirus viral junction region promoter, an alphavirus polymerase recognition sequence and a 3xe2x80x2 polyadenylate tract, wherein said alphavirus RNA, upon packaging into an alphavirus particle and introduction of the particle into a mammalian host cell, has a reduced level of vector-specific RNA synthesis as compared to wild-type alphavirus particle, and the same or greater level of protein encoded by RNA transcribed from the viral junction region promoter, as compared to a wild-type alphavirus particle.\nWithin another embodiment, such RNA vector replicons further comprise a selected heterologous sequence downstream of and operably linked to a viral junction region. Within further aspects of the invention, host cells are provided which contain one of the RNA vector replicons described herein. Within additional aspects of the invention, pharmaceutical compositions are provided comprising RNA vector replicons as described above and a pharmaceutically acceptable carrier or diluent.\nWithin other aspects of the invention, recombinant alphavirus particles are provided, comprising one or more alphavirus structural proteins, a lipid envelope, and an RNA vector replicon as described herein. Within one embodiment, one or more of the alphavirus structural proteins are derived from a different alphavirus than the alphavirus from which the RNA vector replicon was derived. Within other embodiments, the alphavirus structural protein and lipid envelopes are derived from different species. Within further aspects, pharmaceutical compositions are provided comprising a recombinant alphavirus particle as disclosed above and a pharmaceutically acceptable carrier or diluent. Further, mammalian cells infected with such recombinant alphavirus particles are also provided.\nWithin certain embodiments of the invention the above described vectors or particles may further comprise a resistance marker which has been fused, in-frame, with the heterologous sequence. Representative examples of such resistance markers include hygromycin phosphotransferase and neomycin phosphotransferase.\nWithin other aspects of the present invention, methods are provided for selecting alphavirus or recombinant alphavirus vector variants which exhibit the phenotype described herein of reduced, delayed, or, no inhibition of host cell directed macromolecular synthesis. Representative examples of such methods include the use of selectable drug or antigenic markers and are provided in more detail below in Example 2.\nWithin other aspects of the present invention, Togavirus capsid particles are provided which contain substantially no genomic (i.e., wild-type virus genome) or RNA vector replicon nucleic acids. Representative examples of Togaviruses include, for example alphaviruses and rubiviruses (e.g., rubella). Within certain embodiments, the capsid particles further comprise a lipid envelope containing one or more alphavirus glycoproteins. Within other embodiments, the capsid particle further comprises an alphavirus envelope (i.e., the lipid bilayer and the glycoprotein complement). Within related aspects of the present invention, pharmaceutical compositions are provided comprising the above noted capsid particles (with or without a lipid bilayer (e.g., viral envelope containing alphavirus glycoproteins)) along with a pharmaceutically acceptable carrier or diluent. Within further aspects, such capsid particles (with or without a lipid bilayer (e.g., viral envelope containing alphavirus glycoproteins)) or pharmaceutical compositions may be utilized as a vaccinating agent in order to induce an immune response against a desired togavirus.\nWithin further aspects of the invention, inducible promoters are provided comprising a core RNA polymerase promoter sequence, an operably linked nucleic acid sequence that directs the DNA binding of a protein that activates transcription from the core promoter sequence, and an operably linked nucleic acid sequence that directs the DNA binding of a protein that represses transcription from the core promoter sequence. Such promoters may be utilized in the gene delivery vehicles described herein, as well as a wide variety of other vectors known to those skilled in the art.\nWithin other aspects, alphavirus structural protein expression cassettes are provided comprising a 5xe2x80x2 promoter which initiates synthesis of viral RNA from DNA, a nucleic acid molecule which encodes one or more functional alphavirus structural proteins, a selectable marker operably linked to transcription of the expression cassette, and optionally, a 3xe2x80x2 sequence which controls transcription termination. Within one embodiment, such expression cassettes further comprise a 5xe2x80x2 sequence which initiates transcription of alphavirus RNA, a viral junction region promoter, and an alphavirus RNA polymerase recognition sequence. Within another embodiment, the expression cassette further comprises a catalytic ribozyme processing sequence, post-translational transcriptional regulatory elements which facilitate RNA export from the nucleus, and/or elements which permit translation of multicistronic MRNA, selected from the group consisting of Internal Ribosome Entry Site elements, elements promoting ribosomal read through and BiP sequence. Within other embodiments, the selectable marker is operably linked to a 5xe2x80x2 promoter capable of initiating synthesis of alphavirus RNA from cDNA. Within further embodiments the selectable marker is positioned downstream from a junction region promoter and from the nucleic acid molecule which encodes alphavirus structural proteins. Within yet other embodiments, the 5xe2x80x2 promoter is an inducible promoter as described herein. Within another embodiment, the alphavirus structural protein expression cassette further comprises an alphavirus capsid protein gene or other sequence (e.g., a tobacco etch virus or xe2x80x9cTEVxe2x80x9d leader) which is capable of enhancing translation of one or more functional alphavirus structural protein genes located 3xe2x80x2 to the enhancer sequence. Preferably, the capsid protein gene sequence is derived from a different alphavirus than that from which the sequence encoding the alphavirus structural genes is obtained.\nWithin yet other aspects of the invention, alphavirus packaging cell lines are provided comprising a cell containing an alphavirus structural protein expression cassette as described above. In certain embodiments, the alphavirus packaging cell lines are stably transformed with the alphavirus structural protein expression cassettes provided herein. Within related aspects, alphavirus producer cell lines are provided comprising a cell which contains a stably transformed alphavirus structural protein expression cassette, and a vector selected from the group consisting of RNA vector replicons, alphavirus vector constructs and eukaryotic layered vector initiation systems.\nWithin yet other aspects of the present invention, eukaryotic layered vector initiation systems are provided, comprising a 5xe2x80x2 promoter capable of initiating in vivo the 5xe2x80x2 synthesis of RNA from cDNA, a sequence which initiates transcription of alphavirus RNA following the 5xe2x80x2 promoter; a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins, including an isolated nucleic acid molecule as discussed above, an alphavirus RNA polymerase recognition sequence, and a 3xe2x80x2 polyadenylate tract.\nAlso provided are eukaryotic layered vector initiation systems comprising a 5xe2x80x2 promoter capable of initiating in vivo the 5xe2x80x2 synthesis of alphavirus RNA from cDNA, a sequence which initiates transcription of alphavirus RNA following the 5xe2x80x2 promoter, a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins, an alphavirus RNA polymerase recognition sequence, and a 3xe2x80x2 polyadenylate tract, wherein the in vivo synthesized RNA, upon packaging into an alphavirus particle and introduction of the particle into a mammalian host cell, increases the time required to reach 50% inhibition of host-cell directed macromolecular synthesis following expression in mammalian cells, as compared to a wild-type alphavirus particle.\nRelated eukaryotic layered vector initiation system are also provided which comprise a 5xe2x80x2 promoter capable of initiating in vivo the 5xe2x80x2 synthesis of alphavirus RNA from cDNA, a sequence which initiates transcription of alphavirus RNA following the 5xe2x80x2 promoter, a nucleic acid molecule which operably encodes all four alphaviral nonstructural proteins, an alphavirus RNA polymerase recognition sequence, and a 3xe2x80x2 polyadenylate tract, wherein said in vivo synthesized RNA, upon packaging into an alphavirus particle and introduction of the particle into a mammalian host cell, has a reduced level of vector-specific RNA synthesis as compared to wild-type alphavirus particle, and the same or greater level of protein encoded by RNA transcribed from the viral junction region promoter, as compared to a wild-type alphavirus particle.\nRepresentative examples of suitable 5xe2x80x2 promoters for eukaryotic layered vector initiation systems include RNA polymerase I promoters, RNA polymerase II promoters, RNA polymerase III promoters, the HSV-TK promoter, RSV promoter, tetracycline inducible promoter, MoMLV promoter, a SV40 promoter and a CMV promoter. Within preferred embodiments, the 5xe2x80x2 promoter is an inducible promoter as described herein.\nWithin certain embodiments, eukaryotic layered vector initiation systems are provided which further comprise a heterologous sequence operably linked to a viral junction region, and/or a post-transcriptional regulatory element which facilitates RNA export from the nucleus. Within further embodiments, the eukaryotic layered vector initiation systems provided herein may further comprise a transcription termination signal.\nWithin related aspects, the present invention also provides host cells (e.g., vertebrate or insect) containing a stably transformed eukaryotic layered vector initiation system as described above. Within further aspects of the present invention, methods for delivering a selected heterologous sequence to a vertebrate or insect are provided, comprising the step of administering to a vertebrate or insect an alphavirus vector construct, RNA vector replicon, recombinant alphavirus particle, or a eukaryotic layered vector initiation system as described herein. Within certain embodiments, the alphavirus vector construct, RNA vector replicon, recombinant alphavirus particle or eukaryotic layered vector initiation system is administered to cells of the vertebrate ex vivo, followed by administration of the vector or particle-containing cells to a warm-blooded animal.\nWithin other aspects, pharmaceutical compositions are provided comprising a eukaryotic layered vector initiation system as discussed above, and a pharmaceutically acceptable carrier or diluent. Within certain embodiments, the pharmaceutical composition is provided as a liposomal formulation.\nWithin further aspects, methods of making recombinant alphavirus particles are provided, comprising the steps of (a) introducing a vector such as a eukaryotic layered vector initiation system, RNA vector replicon, or alphavirus vector particle as described above into a population of packaging cells under conditions and for a time sufficient to permit production of recombinant alphavirus particles, and (b) harvesting recombinant alphavirus particles. Within related aspects, methods of making a selected protein are provided, comprising the steps of (a) introducing a vector which encodes a selected heterologous protein, such as a eukaryotic layered vector initiation system, RNA vector replicon or alphavirus vector particle described above, into a population of packaging cells, or other cells under conditions and for a time sufficient to permit production of the selected protein, and (b) harvesting protein produced by the vector containing cells. Within yet other aspects, methods of making a selected protein are provided, comprising the step of introducing a eukaryotic layered vector initiation system which is capable of producing a selected heterologous protein into a host cell, under conditions and for a time sufficient to permit expression of the selected protein. Within further aspects, host cell lines are provided which contain a RNA vector replicon as described herein.\nWithin yet other aspects of the present invention, alphavirus vaccines are provided, comprising one of the above-described alphavirus vector constructs, RNA vector replicons, eukaryotic vector initiation systems, or recombinant alphavirus particles, which may or may not express one of the heterologous sequences provided herein (e.g., they may be utilized solely as a vaccine for treating or preventing alphaviral diseases). For example, within one embodiment of the invention, recombinant togavirus particles are provided which have substantially no nucleic acid or RNA vector replicon nucleic acid. Within a further embodiment, recombinant togavirus particles are provided which contain heterologous viral nucleic acids (ie., from a different virus than the togavirus particle). Within yet another embodiment, the recombinant togavirus particle is T=3 or greater.\nWithin further aspects of the invention, recombinant chimeric togavirus particles (either empty, or containing nucleic acids) are provided wherein the viral particle has viral structural components obtained or derived from different Togaviridae (e.g., the capsid protein and glycoprotein is obtained from different alphavirus sources).\nWithin other aspects of the invention, methods for stimulating an immune response within a vertebrate are provided, comprising the step of administering to a vertebrate an alphavirus vector construct, an alphavirus RNA vector replicon according, a recombinant alphavirus particle, or a eukaryotic layered vector initiation system, wherein the alphavirus vector construct, RNA vector replicon, particle, or eukaryotic layered vector initiation system expresses an antigen which stimulates an immune response within said vertebrate (see, e.g., U.S. Ser. No. 08/404, 796 for suitable antigens). Within related aspects, methods are provided for inhibiting a pathogenic agent within a vertebrate, comprising the step of administering to a vertebrate an alphavirus vector construct, an alphavirus RNA vector replicon, a recombinant alphavirus particle, or a eukaryotic layered vector initiation system according, wherein said alphavirus vector construct, RNA vector replicon, particle, or eukaryotic layered vector initiation system expresses an palliative which is capable of inhibiting a pathogenic agent (see. e.g., U.S. Ser. No. 08/404, 796 for suitable palliatives).\nThese and other aspects and embodiments of the invention will become evident upon reference to the following detailed description and attached figures. In addition, various references are set forth herein that describe in more detail certain procedures or compositions (e.g. plasmids, sequences, etc.), and are therefore incorporated by reference in their entirety as if each were individually noted for incorporation."}
{"text": "The described subject matter relates generally to gas turbine engines and more particularly, to arrangements for separating hot and cold flows in gas turbine engines.\nCompact engines require closer packing of components, which in turn requires more crossing of hot and cold gas flows. Without adequate thermal protection, seals, and insulation between these flows, smaller engines suffer from a loss of efficiency. One system developed for certain engines is the mid-turbine frame (MTF), also known as the turbine center frame (TCF) or interturbine frame. This can be disposed between intermediate stages of the turbine section and can have numerous components serving a variety of functions, including as bearing support, engine backbone, combustion gas flow path, coolant flow path, among others.\nDespite tight sealing of the various flows, efficiency gains expected from increased combustion temperatures are not always realized due in part to hot section components being disposed in close proximity to cold section components. As engines are made smaller, the clearances between components, including in the mid-turbine frame area also become tighter. Additional measures must therefore be taken to maintain flow separation and account for higher thermal gradients."}
{"text": "Zipper is one of the most practical inventions of last century, and is extensively used in a wide range of trades and industries. Traditional sliders of zippers may be categorized into automatic locking and non-automatic locking type. An automatic locking slider usually includes a slider body, a pull tab, a locking lever, a plate spring, and a cover. The cover may be fixed to an upper wing of the slider, and attach the pull tab to the slider. The plate spring and the locking lever may be disposed within the cover. The locking lever may receive pressing force generated by the plate spring. Under the force, a locking prong of the locking lever may pass a locking opening on the upper wing to be disposed in the space between two zipper teeth to lock the slider in a position. The pull tab may lift the locking leer when it is pulled. The locking prong may be subsequently removed from between zipper teeth to release the locking of the slider to allow the movement of the slider on zipper. Traditionally, the locking lever, the plate spring and the cover are manufactured by press riveting. The slider body and the cover containing the plate spring and locking lever are also manufactured by press riveting. The manufacturing of a slider is complex and the efficiency is low due to the press riveting and the riveting force is not enough high. In addition, a cover riveting apparatus is required for the press riveting, which increases the production costs of a zipper slider. Even more, there exists the risk that the riveted cover falls off due to inadequate riveting. Tests have shown that the composite pulling force of the slider body and the pull tab of the zipper slider is not stable.\nThe disclosed slider is directed at solving one or more problems set forth above and other problems."}
{"text": "1. Field\nThis disclosure is generally concerned with optical sciences and specifically concerned with mid-IR lenses having high numerical apertures and simultaneously short focal lengths, and further with coupling of such lenses to specialized optical sources.\n2. Related Technology\nLasers are optical sources which provide light of remarkable characteristics. Coherence, high intensity, narrow bandwidth, among others are characteristics common to some lasers which permit their use in conjunction with complex systems to achieve remarkable results. Of particular interest is the highly collimated narrow beam produced by some lasers. Gas lasers have optical resonators defined by two mirrors arranged to produce very narrow beams of planar wavefronts (high spatial coherence). The optical output from such arrangements is highly useful in a great many systems partly because of the nature of the output beam.\nIn contrast to gas lasers, laser diodes are optical sources which convert electrical current to photonic output at a semiconductor PN junction which operates as a laser gain medium. As the geometry associated with such PN junction is generally planar, diodes arranged as lasers sometimes have an asymmetric stripe cross section resulting in output beams having a different degree of divergence in two orthogonal directions. Laser diodes have been used in conjunction with complex optics to condition an output beam to improve its symmetry.\nIn many optical systems, it is desirable to have an axially symmetric collimated beam of planar wavefronts. Yet it can be difficult to arrive at such output from some lasers; especially where a gain medium is constrained by certain geometric features. In those cases, a laser output beam may have less than ideal symmetry and beam shape. Whereas some gas laser systems have an output beam shape which depends on the resonator configuration, diode laser systems sometimes have an output beam conditioned by optics external to the resonator cavity. In some cases, special cooperation between a laser and its output optics can improve system performance.\nLaser diode output beams have been coupled to external systems via micro optical elements. This is certainly true in the case where diode lasers are coupled to fiber optic systems. In such systems the micro optics used may be made from rather conventional materials since laser diode spectra, with very few exceptions, include only those wavelengths suitable for transmission by conventional optics. The art is crowded with semiconductor lasers coupled to external systems via micro optics elements.\nHighly specialized structures recently being used as laser gain media include those known as quantum well stack or sometimes as “quantum cascade”, QC or QCL, laser systems. These structures are semiconductor devices but are strictly not diodes. Rather, a unipolar structure of many layers provide for selective transitions between allowed energy states. When formed in accordance with prescribed designs and coupled with an appropriate optical resonator, a quantum well stack makes an excellent gain medium having good laser performance. In particular, quantum well gain media based laser devices may be arranged with exceptionally wideband gain profiles in the highly useful mid-IR spectrum. In some versions, these lasers are highly tunable and have exceptionally high power output.\nHowever, as in the case with a diode laser beam produced in a quantum well gain medium device may have less than perfectly ideal shape and beam characteristics. The geometry of the gain medium sometimes imparts unwanted influence on the output beam shape. A quantum well stack gain medium is generally formed as a thin stripe semiconductor at the core of a waveguide. The waveguide supports resonant modes which help shape the overall output beam. In particular, a quantum well stack may include an emission facet from which an optical beam leaves the gain medium in a direction substantially normal to the facet surface. The long wavelength or “mid-IR” output beam from a quantum well based laser is typically highly divergent and can be somewhat asymmetric about orthogonal directions. Further, these devices have been constructed by experts in a fashion whereby a sufficiently wide waveguide will support a plurality of transverse oscillation modes. Transverse modes have a greater portion of optical energy ‘off-axis’. This off-axis energy can be lost when a laser is coupled to inefficient systems which suffer aperture clipping effects.\nExperts have now made many systems, instruments, experiments, based upon a quantum cascade lasers and in each and every case without exception, the output of the laser is handled in a fashion whereby the beam is conditioned for use and coupling with other subsystems and experiment or test components. Most generally, the output of a quantum cascade laser is coupled to an experiment by a collimation lens or a parabolic mirror. A quite common way to couple QCLs to external systems is via a mirror having a 50 millimeter focal length and a 50 millimeter clear aperture. While operable, this arrangement suffers from significant edge clipping and the loss is merely tolerated. Since these systems emit mid-IR wavelengths, special lenses which transmit these wavelengths are necessary. Materials such as zinc-selenide, germanium, chalcogonide, or other mid-IR transparent material are used to form these lenses.\nAs the infrared related optical sciences are well developed and quite mature, a great body of lens technologies exists in this area. Particularly, IR lenses configured for imaging tasks are plentiful. With only a cursory effort, an Internet search will generate hundreds of manufactures and designers for lenses useful in IR imaging systems. Imaging lenses are generally compound lenses made up of several individual pieces which function together to perform a high performance imaging relationship with an image plane. While an imaging lens can be used for laser collimation, this arrangement is certainly less than ideal.\nIn addition to IR imaging lens systems, lenses may be arranged for non-imaging or general purpose optical beam handling. Sometimes these lenses are available as “singlets” or single element pieces. In one example of interest, laser collimation applications might call for a plano-convex lens operable at mid-IR wavelengths. Indeed, these lenses are widely available and may be purchased from suppliers of IR lenses. However, laser collimation singlets are designed for and are most suitable for laboratory use with optical beams having large cross-section; that is, laser collimation singlets are most generally large diameter devices of approximately 1 inch. These lenses designed for use in conjunction with special precision optical fixtures designed specifically to couple with low vibration optical benches and permit adjustable alignment and advanced stability features. These lenses typically have a clear aperture of at least 20 millimeters that may be designed with various focal lengths, but typically not with a focal length less than about 20 millimeters. The ratio of focal length to aperture is sometimes called “f#” or ‘f number’. In rare cases, IR singlets having f-number as low as 1 can be found but it has been exceedingly difficult and complex to achieve f-numbers less than 1. These are certainly not available on commercial markets. This is particularly true in the case of longer wavelengths of IR. While near IR systems enjoy the possibility of use of many lens materials, mid-IR systems are highly restrictive in this regard. Materials which transmit mid-IR wavelengths further complicate the possibility of manufacture of micro-optics having low f-number.\nAs many existing quantum cascade laser systems demand a collimated beam, a mid-IR singlet lens is placed in relation to the quantum cascade lasers such that the emission surface is at or near the focal point of the lens. Lenses which transmit mid-IR light are readily available having a focal length approximately 1 inch and an aperture or diameter of about 1 inch. This is the most commonly used optic for coupling a quantum cascade laser to related subsystems.\nWhen a zinc selenide lens of diameter=1″, and focal length=1″ is used to collimate the output of a quantum cascade laser, the lens is placed with its focal plane at the emission aperture of the laser. Because the divergence of the laser output is greater than the numerical aperture of the lens, some light is lost at the lens extremities due to aperture clipping. While the amount of light lost may be as high as 30% of the laser output, this has been tolerated in published systems. This can be even more problematic for lasers which operate with transverse modes. In systems running with excited transverse modes, a greater percentage of optical energy is off-axis and lost at the lens aperture.\nAccordingly, it is very desirable to provide for improved coupling between a mid-IR collimation lens and quantum well gain medium based lasers. First, it is desirable to improve matching of the numerical aperture such that the lens can receive the entire output of the laser. Second, it is highly desirable to reduce the laser energy which is off-axis and susceptible to losses at system apertures (high spatial frequency filtering/clipping). Lasers having a greater portion of energy on-axis can be more effectively coupled to lenses having a numerical aperture less than 1.\nDue partly to materials limits, manufacturing processes available, and industrial experience in IR arts, it has been heretofore impossible to realize mid-IR singlets with numerical aperture greater than about 0.6. For most optical IR systems, this is not a problem as the physical nature of common IR systems does not demand lenses with such performance characteristics. Until the arrival of unipolar quantum well optical sources having characteristic highly divergent mid-IR output beams, and further the arrival of new applications in view of this particular source, there has been little or no effort to discover new devices characterized as high numerical aperture, short focal length, mid-IR lenses, and any such efforts have failed; and indeed left a great and unfulfilled need. This is clearly evidenced as skilled practitioners regularly choose inferior output couplings—i.e. those available.\nExamples of infrared lenses available commercially include those described in the literature as follows:\nNewport Corporation catalog pages including IR lenses describes AMTIR-1 low cost NIR to Mid-IR lenses arranged as singlets in 1″ configurations. These lenses have a low NA or F/# no better than 1.\nComputar Varifocal lenses have short focal lengths but these lenses are arranged as multi-element compound lens systems for IR imaging.\nISP optics catalog pages include ZnSe IR singlets characteristic of the state of the art for IR optics useful in the mid-IR spectrum. These optical elements are as small as 0.5 inch, however, in all cases their f/# is quite large and in no case are these lenses provided for short focal length applications. The lenses are all configured as ‘thin’ lenses having large working distances (and large focal lengths) which contribute to complex optical set-ups suitable for laboratory use only.\nWysokski et al present one of a great many systems representative in the art of a laboratory set-up on optical benches having large optical elements held in specialized optical mounting apparatus. A collimating lens is CL a 1″ diameter f/0.6 lens. This is the typical way that a quantum cascade laser is coupled to an external optical experiment.\nAnother representative illustration of a quantum cascade laser coupled to an experimental set-up in the commonly used configuration is nicely described in the drawing of Hensley's paper: “Recent Updates in QCL-based Sensing Applications” from Physical Sciences, Inc. An optical bench is presented with stabile optical element holders and mounting systems which accommodate 1″ optics common in the industry. The set-up requires a large optical table base for coupling all elements of the optical train which are separated by great distances exposing them to misalignment and vibration problems thereby requiring careful regulation of the apparatus in a guarded environment."}
{"text": "1. Field of Invention\nThe present invention relates to human's mind activities, and more particularly to a mind stimulator, which is capable of substantially cooling down a user's head so as to effectively increase the user's mind activities without hazarding the user's health.\n2. Description of Related Arts\nMany people have an experience that they keep yawning and almost fall asleep during driving. In order to prevent people falling asleep during driving, they either drink coffee as a stimulant or wind down the window to the wind blow into the vehicle. However, such methods are not the best solution to solve the above problem since the coffee contains caffeine that makes people uncomfortable, such as speeding up the rate of heartbeat or causing the hand shaking. In addition, it is a hassle for the passengers when the windows are opened.\nIn order to solve the problem that causes the people sleepy, scientists have discovered that when people are in a drowsy condition, the body temperature, especially the upper head portion, will be raised. Once the body temperature is raised, the nervous system of the human's brain will start slowing down the response so that people will function sluggishly. Therefore, most drivers would like to park their vehicle aside and wash their faces to cool down the body temperature when they are sleepy. In other words, when the body temperature increases, the working efficiency will be reduced.\nFurthermore, the regulating system of the human being cannot automatically cool down the body temperature while the body temperature is raised again. Therefore, people will tend to return to the drowsy condition again such that they must repeatedly wash their faces to keep them awake."}
{"text": "1. Technical Field\nThe present invention relates to methods of manufacturing Group III nitride crystal in which the occurrence of cracking has been reduced.\n2. Description of the Related Art\nGroup III nitride crystal, which is ideally suited to various semiconductor devices including optoelectronic and electronic devices, is finding large-scale utilization. In the manufacture of such III-nitride crystal, various evaporation-deposition techniques—including hydride vapor-phase epitaxy (HVPE), metalorganic chemical vapor deposition (MOCVD), and sublimation—and various liquid-phase deposition techniques—including growth from solution and flux methods—are employed.\nTherein, in order to obtain low-dislocation-density III-nitride crystal in bulk, the epitaxial growth, by an abovementioned technique, onto the principal face of a III-nitride substrate of III-nitride crystal whose constituent-atom type and ratios (that is, chemical composition) are the same as that of the III-nitride substrate, and either the slicing off the III-nitride crystal portion or the removal of the III-nitride substrate by grinding or polishing it away to yield a III-nitride crystal are carried out. (Cf., for example, Int'l./Japanese-Nat'l. App. Pub. No. WO99/23693 [Patent Document 1].)\nIf, however, the dopant type and concentration of III-nitride crystal grown onto the principal face of a III-nitride substrate are different from—even though the constituent-atom type and ratios are the same as—those of the III-nitride substrate, during post-crystal-growth mechanical processing when the III-nitride crystal portion is sliced off or the III-nitride substrate is removed by grinding or polishing it away, the III-nitride crystal is liable to become cracked. This is believed to be due to strain arising between the III-nitride substrate and the III-nitride crystal, owing to the thermal expansion coefficient and lattice constant in III-nitride substrate and in III-nitride crystal being different.\nFurthermore, liquid-phase etching employing a strong acid such as phosphoric acid, or a strong base such as KOH, is ill-suited to III-nitride crystal, inasmuch as the etching speed is slow and the process selectivity is poor."}
{"text": "The present invention is directed to a condensate drain pan for use under the heat exchange coil of an HVAC or refrigeration system. More specifically, the condensate drain pan is multi-functional in that it includes a bottom outlet and a side outlet and in that either outlet can be used on either end of an air conditioning cabinet.\nIn recent years, the air conditioning industry has had a heightened awareness of indoor air quality, particularly relating to the potential health concerns deriving from microbial growth in standing condensate. This standing condensate remains in the drain pans during the air conditioning unit's off cycle, and may be subsequently transported to an occupied space by way of the indoor air stream. For this reason, the air conditioning industry has recognized the need for a periodic cleaning of drain pans. Drain pans are positioned under a condensing heat exchange coil to collect the condensate as that condensate condenses and drips off the coil.\nTypical past industry practice has provided a flat condensate pan with condensate connections at only one end of the drain pan. In installing such drain pans, particularly in retrofit applications, it is not uncommon for the air conditioning unit to be improperly leveled during installation. If an air conditioning unit is improperly leveled such that the slope of the unit is away from the single drain pan connection, condensate will pool at the end of the drain pan opposite the drain connection. In more extreme cases, condensate will overflow the drain pan resulting in unwanted water in the base of the air conditioning unit and, potentially, in the occupied space and structure. Past solutions to this situation have included (a) removal of the air conditioning unit and re-leveling of the mounting curb at great expense, or (b) if the unit is large enough, providing a connection at either end of the drain pan. In many cases, the problem is simply ignored due to the difficulty in correcting the problem.\nFurthermore, access for cleaning of drain pans in the past has typically been through a time consuming removal of the exterior panels of the air conditioning units, followed by an attempt to clean the drain pan through a limited access aperture. As a result, the drain pan was often not cleaned, or was improperly cleaned. In some cases, chemical \"pills\" to counteract microbial growth have been added to drain pans in lieu of full drainage and proper cleaning.\nIndustry standards have been under development that make it unacceptable to provide air conditioning unit designs that allow condensate to pool in drain pans. The practical industry response to this has been to provide drain pans with sloped surfaces to ensure complete drainage. However, this assumes that the air conditioning units are properly leveled. With a sloped drain pan there is no longer a way to provide connections on both ends of the drain pan without one end of the drain pan being dysfunctional, thereby eliminating the second solution above and leaving only the solution of incurring the expense of re-leveling the air conditioning unit if the unit is improperly installed.\nRegarding bottom condensate connections, it is becoming common practice in the building/roofing industry to eliminate as many roof penetration apertures as possible due to the substantial warranties and potential liabilities associated with roof leaks in that industry. As such, there is a growing preference to bring the condensate disposal piping of the air conditioning unit from the space up through the bottom of the unit to eliminate a roof penetration and a potential leak related liability. Solutions to date have been (a) to route condensate piping out of the side of the unit and back in through the side of the rooftop curb, or (b) to provide both bottom and side outlet connections, requiring whichever outlet connection is not in use to be plugged or capped with a potential for leaks, or (c) to mount the drain pan in an elevated position within an air conditioning unit such that the piping could be routed out either the side or through the base of the unit.\nIt would be desirable to provide a drain pan which is easily cleanable, which is sloped, which can be installed from either side, and which includes both a bottom outlet and a side outlet."}
{"text": "1. Field of the Invention\nThis invention relates to pressurized containers and to methods for pressurizing and filling them. In accordance with a first aspect of the invention, the containers are completed by the container manufacturer and shipped ready to be filled. In a preferred embodiment a propellant is introduced into the completed container by the manufacturer before the container is shipped to a filler to be filled with product. According to a second aspect of the invention the container is pressurized and filled in a way to ensure that the container is not excessively pressurized during filling and an adequate pressure is maintained in the container until all or substantially all of the product is depleted during use.\n2. Prior Art\nPressurized containers are used to dispense a variety of products, including paint, lubricants, cleaning products, food items, personal care products such as hair spray, and the like. Pressure for dispensing these products is provided by a propellant placed in the container. In some prior art systems the product and propellant are stored separately in the container, i.e., separated by a barrier, e.g. a piston or bag, commonly referred as a barrier pack system. In other systems the product and propellant are stored together under pressure in the container. Dispensing of the product occurs when a discharge valve or nozzle is opened, permitting the pressurized product to be forced out through the nozzle, usually as a spray, stream, or foam. As product is depleted from the container, the pressure exerted by the propellant decreases, especially evident when compressed gases are used as the propellant, and the propellant pressure may become diminished to the extent that all of the product cannot be dispensed from the container, or a desired characteristic, e.g., atomization, is not achieved.\nIn addition to the propellant component, many products, e.g., hair spray, require a carrier, e.g., alcohol, or combinations of alcohol with water or other volatile solvents that dry quickly upon discharge from the container. Other volatile solvents or propellants that can be used in these systems include volatile organic compounds (VOCs) such as propane, isobutane, dimethyl ether, and the like, but their use is limited due to environmental concerns. For instance, under some current regulations no more than 55% of the contents of the container can comprise a VOC. In an aerosol dispenser, as much as 25% of the VOC could be required for use as a propellant, leaving about 30% VOC in the product. The balance of the product would be the active ingredients and water, which does not dry as quickly as the VOC, resulting in a “wet” product when used.\nCarbon dioxide (CO2) is useful as an aerosol propellant, but its use has been limited due to the fact that it is normally placed in the container as a pressurized or compressed gas, and in conventional systems the drop-off in pressure is excessive as the product is depleted and the volume occupied by the propellant increases. For example, in a typical situation the starting pressure might be 90-125 psig and the finishing pressure only 20 or 30 psig.\nConventional barrier pack systems typically comprise a can made of aluminum, steel, plastic, or other suitable material, with a barrier in the can between the product and the propellant. The barrier normally comprises a piston reciprocable in the can, or a collapsible bag in which the product is contained. In accordance with conventional practice, barrier pack cans are shipped empty from the manufacturer to a location where the can is to be filled, either with a piston in place in the can or a bag attached to the valve or the dome closing the end of the can. The filler adds the product, crimps and seals the valve in place in the opening provided for that purpose in the domed top of the can, and then injects the propellant.\nIf the barrier pack is of the type having a piston, the filler normally introduces product, e.g., a gel, through the opening in the domed top and into the can above the piston. The aerosol valve is then fitted and sealed to the can, and a propellant such as, e.g., isobutane, a VOC, is introduced under a predetermined pressure into the can beneath the piston through a sealing plug in the bottom of the can. If a liquefied propellant is used, some of it vaporizes until an equilibrium pressure is reached. The pressurizing propellant forces the piston up, placing pressure on the product so that it is discharged through the valve when the valve is opened.\nIn barrier packs utilizing a bag wherein the bag is affixed to the valve body on the bottom side of the valve cup with an undercup gasser, the filler introduces a propellant around the valve and into the can outside the bag, crimps the can, and then introduces product into the bag through the valve. Alternatively, a second method utilizes a plastic bag that is pre-inserted into the can and that has a formed one-inch neck shaped to fit the curl of the can, which allows product to be filled before the valve is applied and sealed. Propellant is then added through the sealing plug in the bottom of the can. The propellant exerts pressure on the bag, forcing product out through the valve when the valve is opened.\nIn those conventional systems wherein the propellant is mixed in the container with the product, the can manufacturer ships an empty container to the filler, who then places a desired quantity of product into the container, attaches and seals the valve, and then injects propellant through the valve to pressurize the product.\nIn order to promote dissolution of the propellant into the product as the propellant is being introduced, some prior systems shake the container, thereby reducing the pressure spike or over-pressurization that occurs when the propellant is first charged into the container and thus avoiding deformation of the can. However, these prior art systems have not been entirely satisfactory because of slower gassing and the shaking required.\nVarious other systems have been developed in the prior art for storing a reserve supply of propellant and adding it to the container as product is depleted, so that propellant pressure is maintained at a desirable level until a suitable amount of the product is dispensed from the container. Examples of such systems are described in applicant's prior issued U.S. Pat. Nos. 6,708,844 and 7,185,786, and applicant's prior copending U.S. application Ser. No. 11/250,235, filed Oct. 14, 2005, all of which are incorporated in full herein by reference.\nCommon to the foregoing systems is the need for the filler to provide machinery for completing manufacture and/or assembly of the final product, and in the case of pressurized aerosol dispensers to inventory propellants and solvents in addition to the product. For many small fillers, in particular, this is a burdensome requirement due to the cost of the necessary machinery to complete manufacture of the containers and to store propellant gases, and when applicable the cost of carrying insurance and maintaining appropriate storage facilities for required propellants and solvents.\nIt would be advantageous to have an economical, efficient, and environmentally safe system and method for filling and pressurizing containers, wherein completed containers are shipped by the container manufacturer to the filler so that the filler does not require the necessary equipment to complete the vacuum crimping, propellant gas injection, gas storage tanks, and pumping equipment to complete the manufacture of the pressurized product, and does not need to incur the cost of carrying insurance and maintaining manufacturing and appropriate storage facilities for required propellants. Moreover, it would be advantageous to have a system and method for filling and pressurizing containers wherein the initial starting pressure is not excessive and satisfactory pressure is maintained throughout the useful life of the container."}
{"text": "The current methods for preparing florpyrauxifen-benzyl (benzyl 4-amino-3-chloro-5-fluoro-6-(4-chloro-2-fluoro-3-methoxyphenyl)picolinate; I) are described in U.S. Pat. Nos. 8,609,855 and 8,871,943.\n"}
{"text": "1.\nAs of the Invention\nThe invention relates to a flow line monitor pilot valve assembly which automatically activates a spool valve assembly to block and bleed control fluid for manipulation of a downstream valve mechanism to closed position when abnormal pressures in a monitored fluid are detected. The spool valve is manually or mechanically manipulated at any time to open or close the downstream valve.\n2. Description of the Prior Art As a safety precaution, oil and gas wells are required to contain one or more subsurface safety valves, which are typically of ball-type construction and are manipulated between open and closed position by variation of hydraulic control pressure. The safety valves are manipulatable to closed position when detection devices immediate the well respond to an unfavorable condition, such as an increase in temperature, indicating a fire, or by abnormal increase or decrease in well production flow, indicating a flow line plug or rupture. The monitoring devices transmit the abnormal signal to a device which causes automatic loss of hydraulic control fluid to the safety valve, and the safety valve is manipulated to closed position.\nPilot relays heretofore utilized have incorporated pressure unbalanced poppet-type valving mechanisms because spool valves having very close tolerance metal sleeves result in high friction being encountered across the valve as the head is shifted with respect to its seat. Thus, it would be desirable to have a valving mechanism which does not encounter friction as the head is moved relative to the seat.\nSome prior art pilot relays are pressure unbalanced and require pneumatic signal to allow hydraulic pressure to the downhole safety valve. In these systems, the pneumatic source is typically tapped from the well, so that when the well is shut in, there is no source of control fluid to reopen the valves. This, in turn, has required utilization of mechanical defeat mechanisms or hand pumps to activate the system to manipulate the valve again to open position. These start-up techniques, if not mechanically released, will defeat the entire safety system.\nValving mechanisms heretofore known and utilized have required the incorporation of numerous valves and complicated circuitry. The multiplicity of valves has been utilized to interface many sensing signals, usually at a minimum of one valve per signal. In these systems, manual valves are required to by-pass sensing signals for start-up. Additionally, other manual valves are required to shut down the system. Frequently, hand pumps or mechanical defeat mechanisms have been required to start up the systems. Not only are these valving systems complicated, but they are expensive, result in frequent service problems, and, because of their complexity, have a higher probable failure rate, and are thus unreliable.\nSome prior art valving means have utilized spherical elements which essentially act as a valve head which is selectively movable onto and off of an adjacent seat by means of a spool element which is shuttled to manipulate a spherical element. For example, U.S. Pat. No. 2,574,335 discloses two sets of two ball elements each which are opened and closed at the same time in a fluid pressure motor device for control systems. The shuttle device contains a conical carrier surface manipulating a first ball on the conical or beveled surface to move an enlarged ball thereon between open and closed position within flow passageways.\nU.S. Pat. No. 2,891,518 utilizes spherical or ball elements which are not normally free traveling or sealingly engagable upon a seat by differential pressure caused by fluid flow.\nU.S. Pat. No. 2,967,544 also is typical of the prior art utilizing spherical elements as a valve head. This patent utilizes compressed spring elements and a mechanical manipulator to shift the balls between opened and closed position with respect to their ports. The spherical or ball elements are not carried by a spool which is shuttleable.\nU.S. Pat. No. 3,007,492 discloses a pilot valve for fluid flow systems utilizing a spherical or ball element as the valve head. The ball elements are contained within a shuttleable spool element which is only shuttleable by exertion of a spring element in one direction. A piston head in a piston chamber is utilized to shuttle the spool in the opposite direction against the spring.\nU.S. Pat. No. 3,053,279 discloses a directional valve containing a plurality of ball elements manipulatable by a spool to control flow from a second passageway to a first passageway. In one position of the spool, all of the balls are positionable on each of the respective seats. This valve is utilized to control flow from a second passageway to a first passageway, and flow therebetween is controlled by the positioning of a ball within the first passageway.\nThe present invention obviates the problems set forth above and differs substantially from the prior art. The present invention allows shifting of the spool element in response to detected abnormal pressure of a monitored fluid to block and bleed control fluid without affecting manual control capabilities. Additionally, the present invention does not require by-pass valves, mechanical defeat devices, or hand pumps for system start-up. Additionally, the present invention does not rely upon pressure unbalanced poppet-type valve mechanisms, but provides a bi-stable valve which tolerates exposure to high pressures. Additionally, the present invention permits easy manual or mechanical shuttling of the spool assembly at any time. Moreover, the present invention differs substantially from the prior art by utilization of a spool element containing spherical or ball means carried thereby which are selectively engageable onto and wipable off of their respective seats.\nOther advantages of the present invention will be apparent from a reading of the FIGS., the specification, and the claims below."}
{"text": "Transmission Control Protocol/Internet Protocol (TCP/IP) is a network communication protocol in which a sender transmits a packet to a receiver, sets a timer in association with the transmission, and waits until a timeout for the receiver's acknowledgement. If the acknowledgement is not received before the timeout, or if an acknowledgement is received but indicates error, the sender retransmits the packet to the receiver.\nEnd-to-end paths between sender and receiver in TCP/IP communications can be a sequence of links, and any of these links can drop packets. Regardless of where the drop occurs though, retransmission in an end-to-end TCP/IP communication can consume bandwidth in every link, from the sender to the receiver.\nHigh packet drop can be exhibited, in some network applications, in the last link or loop to the receiver, or “last mile.” There can be various reasons. For example, channel conditions can be degraded due to weather, radio interference, and other time-varying phenomena. In addition, upgrading to reduce packet loss in the last mile can be costly."}
{"text": "A number of different applications of catadioptric imaging systems for far-field and near-field interferometric confocal and non-confocal microscopy have been described such as in commonly owned U.S. Pat. No. 6,552,852 entitled “Catoptric And Catadioptric Imaging Systems” and U.S. Pat. No. 6,717,736 entitled “Catoptric And Catadioptric Imaging Systems;” U.S. Provisional Patent Applications No. 60/447,254, filed Feb. 13, 2003, entitled “Transverse Differential Interferometric Confocal Microscopy,”; No. 60/448,360, filed Feb. 19, 2003, entitled “Longitudinal Differential Interferometric Confocal Microscopy for Surface Profiling,”; No. 60/448,250, filed Feb. 19, 2003, entitled “Method and Apparatus for Dark Field Interferometric Confocal Microscopy,”; No. 60/442,982, filed Jan. 28, 2003, entitled “Interferometric Confocal Microscopy Incorporating Pinhole Array Beam-Splitter,”; No. 60/459,425, filed Apr. 1, 2003, entitled “Apparatus and Method for Joint Measurement Of Fields Of Scattered/Reflected Orthogonally Polarized Beams By An Object In Interferometry,”; No. 60/485,507, filed Jul. 7, 2003, entitled “Apparatus And Method For High Speed Scan For Sub-Wavelength Defects And Artifacts In Semiconductor Metrology,”; No. 60/485,255, filed Jul. 7, 2003, entitled “Apparatus and Method for Ellipsometric Measurements with High Spatial Resolution,”; No. 60/501,666, filed Sep. 10, 2003, entitled “Catoptric And Catadioptric Imaging Systems With Adaptive Catoptric Surfaces,”; No. 60/602,046, filed Aug. 16, 2004, entitled “Apparatus And Method For Joint And Time Delayed Measurements Of Components Of Conjugated Quadratures Of Fields Of Reflected/Scattered Beams By An Object In Interferometry,”; and U.S. patent applications Ser. No. 10/778,371, filed Feb. 13, 2004, entitled “Transverse Differential Interferometric Confocal Microscopy,”; Ser. No. 10/782,057, filed Feb. 19, 2004, entitled “Longitudinal Differential Interferometric Confocal Microscopy for Surface Profiling,”; Ser. No. 10/782,058, filed Feb. 19, 2004, entitled “Method and Apparatus for Dark Field Interferometric Confocal Microscopy,”; Ser. No. 10/765,229, filed Jan. 27, 2004, entitled “Interferometric Confocal Microscopy Incorporating Pinhole Array Beam-Splitter,”; Ser. No. 10/816,180, filed Apr. 1, 2004, entitled “Apparatus and Method for Joint Measurement Of Fields Of Scattered/Reflected or Transmitted Orthogonally Polarized Beams By An Object In Interferometry,”; Ser. No. 10/886,010, filed Jul. 7, 2004, entitled “Apparatus And Method For High Speed Scan For Sub-Wavelength Defects And Artifacts In Semiconductor Metrology,”; Ser. No. 10/886,157, filed Jul. 7, 2004, entitled “Apparatus and Method for Ellipsometric Measurements with High Spatial Resolution,”; and Ser. No. 10/218,408, filed Sep. 10, 2004, entitled “Catoptric And Catadioptric Imaging Systems With Adaptive Catoptric Surfaces,”. In addition, U.S. patent application Ser. No. 10/218,201, entitled “Method for Constructing a Catadioptric Lens System,” filed Apr. 1, 2004 described one way to make some of these catadioptric lens systems. These patents, patent applications, and provisional patent applications are all by Henry A. Hill and the contents of each are incorporated herein in their entirety by reference.\nIn each of the applications of catadioptric imaging systems for each of the cited U.S. patents, U.S. patent applications, and U.S. Provisional Patent Applications, a beam-splitter is incorporated in generating an image of an object with zero optical aberrations for a measurement object located on the optic axis of the imaging system. The beam-splitter is located at an interface between relatively thick optical elements of the catadioptric imaging systems. The optical elements contribute off-axis aberrations and cause a significant portion of optical paths in the catadioptric imaging systems to comprise a transmitting refractive medium such as fused silica or CaF2.\nIn each of the applications of catadioptric imaging systems for each of the cited U.S. patents, U.S. patent applications, and U.S. Provisional Patent Applications, tight tolerances are generally placed on the manufacture of optical elements. In addition to the tolerances normally encountered in designing a diffraction limited imaging system, there are additional tolerances imposed in interferometric confocal and non-confocal microscopy applications. The additional tolerances are for example on surfaces of certain elements with respect to radii of curvature and on relative locations of centers of curvature of the surfaces of the certain elements.\nThe additional tolerances lead to improved performance of a catoptric or a catadioptric imaging system, e.g., with respect to increasing the average intensity of desired images by a factor of approximately 2 and reduced intensity of spurious beams by one or more order of magnitudes, and in addition make it possible to realize interferometric reduction of background fields. The interferometric reduction of background fields leads to a reduction of statistical errors. The increase in intensity of desired images and the reduction of statistical errors lead to an increase in signal-to-noise ratios and to a concomitant increase in throughput of a metrology tool using the catoptric or catadioptric imaging system. The interferometric reduction of background fields further leads to a reduction of systematic errors. A consequence of the reduction of systematic errors is a reduction of the computational task required to invert arrays of measured interference signal values to a multi-dimensional image of a measurement object.\nThe cited U.S. patents, U.S. patent applications, and U.S. Provisional Patent Applications further teach the use of adaptive catoptric surfaces in a catoptric or catadioptric imaging system. The use of adaptive catoptric surfaces in a catoptric or catadioptric imaging system makes it possible to relax tolerances on the surface figures of elements, to relax tolerances on locations of surfaces of the elements in the catoptric or catadioptric imaging system, and to compensate for certain optical aberrations such as may be introduced by the pellicle or aperture-array beam-splitter. The factor by which the tolerances may be relaxed on the surface figures is of the order of 5 for certain of the elements. The use of adaptive catoptric surfaces in a catoptric or catadioptric imaging system further makes it possible to introduce a vertical or lateral scan of a measurement object or substrate being imaged at slew rates higher than possible and/or practical when the vertical or lateral scan must otherwise be introduced either by translations of an entire catoptric or catadioptric imaging system and associated optics and detector systems or translations of the measurement object or substrate, e.g., a 300 mm wafer, and the measurement object or substrate support system.\nThe cited U.S. patents, U.S. patent applications, and U.S. Provisional Patent Applications further teach the replacement of a beam combining beam-splitter in an interferometric imaging system with a thin fluorescent layer or interface.\nThe cited U.S. patents, U.S. patent applications, and U.S. Provisional Patent Applications also teach the use of an N-dimensional bi- and quad-homodyne detection methods."}
{"text": "Ground-mounted photovoltaic solar panels are conventionally supported on solar panel mounting racks. Commercially available solar panel racks are typically produced using aluminum extruded sections or steel roll formed sections in order to provide the structural strength required to withstand loads associated with outside conditions such as wind and snow."}
{"text": "1. Field of the Invention\nThe present invention relates to an engine exhaust gas purification device.\n2. Description of the Related Art\nIf a ratio of the total amount of air fed into the exhaust passage, the combustion chamber and the exhaust passage upstream of a certain position in the exhaust passage to the total amount of fuel fed into the exhaust passage, the combustion chamber and the exhaust passage upstream of the above-mentioned position in the exhaust passage is referred to as an air-fuel ratio of the exhaust gas at that certain position, in a known engine, a lean air-fuel mixture is burned in the combustion chamber, and a three way catalyst is arranged in the exhaust passage, and a NOx absorbent, which absorbs NOx when the air-fuel ratio of the exhaust gas flowing into the NOx absorbent is lean and releases absorbed NOx when the air-fuel ratio of the exhaust gas flowing into the NOx absorbent is rich, is arranged in the exhaust passage downstream of the three way catalyst (see PCT international publication WO93/07363).\nNamely, in an engine, the amount of production of NOx becomes maximum when an air-fuel mixture which is slightly leaner than the stoichiometric air-fuel ratio is burned, and the amount of production of NOx becomes smaller as the air-fuel ratio of air-fuel mixture is made leaner or richer with respect to the air-fuel ratio which is slightly leaner than the stoichiometric air-fuel ratio. In addition, when a lean air-fuel mixture is burned, the amount of unburned HC and CO discharged from the engine becomes considerably small, as compared with the case where a rich air-fuel mixture is burned. Accordingly, in the above-mentioned engine, the air-fuel ratio of the air-fuel mixture to be burned in the combustion chamber is made considerably lean to reduce the amount of NOx and unburned HC and CO discharged from the engine as much as possible. In addition, at this time, the NOx discharged from the engine is caused to be absorbed by the NOx absorbent as much as possible, and the unburned HC and CO discharged from the engine is caused to be oxidized by the three way catalyst as much as possible to reduce the amount of NOx and unburned HC and CO discharged into the outside air.\nAmong the harmful components in the exhaust gas, NOx is a component which is the most difficult to purify and, if this NOx can be reduced by the three way catalyst, it is possible to considerably lower the amount of NOx in the exhaust gas. However, if the air-fuel ratio of the exhaust gas is lean, the three way catalyst does not carry out the reducing operation of NOx. Accordingly, if the air-fuel ratio of air-fuel mixture is made lean, i.e., the air-fuel ratio of the exhaust gas is made lean as in the above-mentioned engine, NOx discharged from the engine is not reduced by the three way catalyst and, accordingly, in the above-mentioned engine, NOx passes through the three way catalyst without being reduced and then flows into the NOx absorbent.\nNamely, in the above-mentioned engine, the first stage of the reducing operation of the amount of NOx is carried out by making the air-fuel ratio of the air-fuel mixture lean. In this case, it is true that the amount of NOx discharged from the engine is reduced due to this first stage of the reducing operation of the amount of NOx, but the absolute amount of NOx discharged from the engine is still large. Then, although NOx discharged from the engine flows into the three way catalyst, the amount of NOx is not reduced in the three way catalyst at all. Then, the NOx flows into the NOx absorbent, and the second stage of the reducing operation of the amount of NOx is carried out therein due to the absorbing operation of NOx by the NOx absorbent.\nIn this way, in this engine, the reducing operation of NOx is carried out in two stages, but the main reducing operation of the amount of NOx relies on the absorbing operation of NOx by the NOx absorbent. However, even if the NOx absorbent is used, particularly when the absolute amount of NOx flowing into the NOx absorbent is large, all the NOx cannot always be absorbed in the NOx absorbent. Accordingly, it is difficult to make the amount of NOx, discharged into the outside air, almost zero as long as the main reducing operation of NOx is made to rely on only the absorbing operation of NOx by the NOx absorbent."}
{"text": "Electronic money cards are now widely used for transacting business. Generally, these electronic money cards are credit card size electronic devices that have a data port that can be connected into an electronic money transaction machine that increases or decreases the balance stored on the money card, depending upon the transaction performed. For example, the electronic money card can be inserted into an electronic bank machine that can increase the balance stored on that card to represent thereby the withdrawal of funds from the bank. Conversely, the electronic money card can be inserted into an electronic cash register machine that can decrease the balance on the card responsive to a purchase of goods.\nElectronic money cards work well for carrying and storing data signals representative of electronic money. However, unlike conventional hard currency, electronic money signals are not readily discernible by a user wishing to determine the amount of money available to that user. Accordingly, electronic money cards require a balance reader that allows the user to determine the balance stored on that card. Typically, these balance readers are electronic circuit card assemblies sized to receive a money card, with a port for slidably receiving the electronic money card and with a display that can display to a user the balance stored in that card. Although these balance readers generally work well, they tend to be bulky, and complicated plastic cases that are cumbersome to carry.\nAccordingly, it is an object of the invention to provide balance readers that are more facile to use and more convenient to carry.\nOther objects of the invention shall be made obvious by the review of the attached figures and upon study of the associated detailed description."}
{"text": "In memories, individual memory cells are aligned along rows and columns. Each row has a word line to which cells along that row are attached. The cells along a row are enabled when the word line to which the cells are attached is enabled. It is important that only the word line for the selected row be enabled. Due to capacitive coupling of various clock signals with fast switching times, it is possible to enable a word line of an unselected row with capacitive coupling unless measures are taken to prevent this. This is typically achieved by clamping deselected word lines to ground which, although effective, has required large pull-down transistors, the layout of which tends to limit the pitch of the word lines, i.e., limit how close together the word lines can be.\nAlso typical of memory row decoders is that they are NOR decoders which have a relatively high input capacitance. This causes significant loading to the address buffers which provide the address signals to the decoders. The address signals are input to the gates of transistors which form the NOR decoder. The output of the NOR decoder is on the drains of these transistors. The output is thus inverted from the input. Consequently, the gate to drain capacitance is always charged between a logic high and a logic low. Consequently, a change in state of the NOR decoder not only requires that the gate to drain capacitance be charged but also be discharged to the opposite polarity. This also adversely affects the speed of the decoder."}
{"text": "1. Technical Field\nThe present disclosure relates to monitoring systems and, particularly, to a monitoring system and a monitoring method.\n2. Description of the Related Art\nConventional monitor systems typically include a number of telecameras, a number of monitors, and a number of video cassette recorders. The telecameras aim at a scene so as to capture images of the scene from different directions, and output surveillance videos to the monitors for monitoring in real time. The surveillance videos are stored in the video cassette recorders.\nOne of the disadvantages of utilizing such a monitoring system is that in order to find a target person in the surveillance videos, all the surveillance videos need to be reloaded from the cassette recorders, and checked one by one. This manual searching task is time-consuming and inefficient.\nWhat is needed, therefore, is a monitoring system and a monitoring method, which can overcome the above-described problem."}
{"text": "Charge-transporting thin-films composed of organic compounds are used as the emissive layer or charge injection layer in organic electroluminescence (organic EL) devices. Processes for forming such thin-films are broadly divided into dry processes such as vapor deposition and wet processes such as spin coating. In comparing dry processes and wet processes, wet processes are capable of efficiently producing thin-films having a high flatness over a large surface area. Therefore, in the organic EL field where thin-films of larger surface area are desired, organic thin-films are often formed by wet processes.\nIn recent years, there has been a steadily growing desire for higher functionality, greater multifunctionality, smaller size, etc. in manufactured products within the organic EL field. Under such circumstances, to achieve, for example, lighter weights and thinner dimensions, substrates composed of organic compounds have come to be used in place of glass substrates, and so the treatment temperatures that can be used in manufacturing processes have become lower than in the past. As a result, there exists today in the organic EL field an increasing desire for varnishes which can be baked at lower temperatures and which, moreover, even in such cases, provide thin-films endowed with good charge transportability."}
{"text": "This invention is related generally to a dual polarized aperture coupled microstrip patch antenna system. A variety of aperture coupled patch antenna systems are available and known. Prior art dual polarization aperture coupled patch antenna systems suffer from either small bandwidth, poor cross polarization port-to-port isolation (about -20 dB) or require complicated and expensive feed structures including multiple feeding network layers or jumpers or crossovers.\nAccordingly, a need arises for a dual polarized aperture coupled patch antenna system having an uncomplicated, inexpensive feeding structure with effective cross polarization port-to-port isolation."}
{"text": "1. Field of the Invention\nThe present invention relates to a printing apparatus and a method of acquiring correction value. Specifically, the invention relates to a technique to acquire a correction value to correct an error in conveying a printing medium used in an inkjet printing apparatus.\n2. Description of the Related Art\nAn inkjet printing apparatus has a print head that has a fine-nozzle array, and ink is ejected from each nozzle in accordance with printing data. The ejected ink forms dots on the printing medium to form an image. Accordingly, to form a high-quality image, it is important that the dots should be formed on the printing medium at intended positions. The displacement of the dot-formation position has to be avoided as much as possible. Some of the various causes of such displacement deviation are: difference in shape amongst the nozzles of the print head; noise factors, such as the vibrations of the apparatus that occur while the printing is being carried out; and the distance between the printing medium and the print head. The inventors of the present invention have discovered that one of the significant causes for such displacement deviation of the dot-formation position is the lack of accuracy in conveying the printing medium. One of the commonly used conveying unit for the printing medium is a roller (a conveying roller). Conveying the printing medium by a desired distance can be achieved by rotation of the conveying roller by a designated angle with the conveying roller being pressed onto the printing medium. Here, the accuracy in the conveying of the printing medium depends, to a significant extent, on the eccentricity of the conveying roller.\nFIGS. 33, 34A and 34B, and 35 illustrate cross sectional shapes of various conveying rollers. The conveying roller of FIG. 33 has a perfectly-circular cross-sectional shape, and has its central axis aligned exactly with its rotational axis. The conveying roller of FIGS. 34A and 34B has a cross-sectional shape that is not a perfect circle. The conveying roller of FIG. 35 has its rotational axis offset from its central axis.\nAssume such a case as shown in FIG. 33, or, to be more specific, a case where the cross-sectional shape of the conveying roller is a perfect circle and where the central axis of the conveying roller is aligned exactly with its rotational axis. In addition, further assume that the rotational angle to convey the printing medium is uniform. Then, every rotation of the conveying roller by an angle R constantly gives a particular length (L0) in the circumferential directions (length of arc). Accordingly, every position within the conveying roller always gives a uniform amount of conveying the printing medium that is conveyed while being in contact with the conveying roller.\nContrasting outcomes are obtained by a conveying roller with an ellipsoidal cross-sectional shape as ones shown in FIGS. 34A and 34B. Such a conveying roller gives different amount of conveying even when the conveying roller rotates by the same angle R. This difference in the amount of conveying depends on the rotational position of the conveying roller. To be more specific, for the rotational position shown in FIG. 34A, the printing medium is conveyed by an amount L1 while for another rotational position shown in FIG. 34B, the printing medium is conveyed by an amount L2. Here, the lengths L0, L1, and L2 has such a relationship as L1>L0>L2. That is to say, a periodical variation in amount of conveying the printing medium occurs, and the variation depends on the period of the conveying roller.\nAlternatively, as in the case of FIG. 35, the offsetting of the rotational axis of the conveying roller from the central axis O that is intended to be the rotational axis may sometimes cause the amount of conveying the printing medium to vary periodically in response to the period of the conveying roller. To be more specific, assume cases where the rotational axis is offset from the central axis O and is positioned at either the point A or the point B shown in FIG. 35. In these cases, the same rotational angle α produces different amounts of conveying. Such difference in conveying amount results in a periodical variation in the conveying of the printing medium. Here, the variation depends on the period of the conveying roller.\nThe eccentricity of the roller, which has been mentioned above, includes these above-described states. Specifically, included are a state where the roller has a cross-sectional shape that is not a perfect circle, and a state where the conveying roller has its rotational axis offset from its central axis. In the case of an ideal accuracy being achieved in conveying, the image should be printed in such a way as shown in the schematic diagram of FIG. 36A. With the above-mentioned eccentricity, however, the printed image will be an uneven image with stripes that appear periodically in the conveying direction as shown in FIG. 36B while the period is the same as the amount of conveying corresponding to a full rotation of the conveying roller.\nThe amount of eccentricity for the conveying roller is usually controlled so as to stay within a certain range. The stricter the standard for the amount of eccentricity is, the lower the yielding of the conveying roller becomes. Accordingly, the printing apparatus thus produced becomes more expensive. For this reason, an excessively strict standard for the amount of eccentricity is not preferable.\nTo address the above-mentioned problem, various measures have been proposed. Different correction values for the conveying errors are set for different phases of the conveying roller so that even an eccentric conveying roller can achieve a steady amount of conveying as similar to the case of a conveying roller with a perfectly-circular cross-sectional shape and with its rotation axis being aligned exactly with its central axis (Japanese Patent Laid-Open No. 2006-240055 and Japanese Patent Laid-Open No. 2006-272957). To be more specific, correction to reduce the amplitude of the fluctuation in amount of conveying with a period equivalent to the circumferential length of the conveying roller can be done by applying a periodic function with the same period and reversed polarity.\nBesides the eccentricity that is mentioned above, variations in outer circumference, or outer diameter, of the roller is another important cause for lowering the accuracy in conveying. With such variations in outer diameter of the roller, rotation of a roller by a rotation angle determined for a roller with a reference outer diameter will not produce a predetermined amount of conveying. Specifically, use of a roller with an outer diameter that is larger than the standard outer diameter produces a larger amount of conveying while use of a roller with an outer diameter that is smaller than the standard outer diameter produces a smaller amount of conveying. Accordingly, even when the amplitude of the variation is reduced by the above-described correction, the range of variations which is maximum and which exceeds a certain amount of conveying error causes unevenness that appears in the image. This means that to achieve the printing of a high-quality image without unevenness requires not only the lowering of the influence of the eccentricity but also the lowering of the influence of the variations in the outer diameter of the conveying roller.\nAn example of the techniques to achieve the printing of a high-quality image reducing unevenness is disclosed in Japanese Patent Laid-Open No. 2002-273956. In the disclosed technique, the correction value to correct the conveying error caused by the variations in the outer diameter of the conveying roller (correction value for outer-diameter) is acquired. Also acquired is the correction value to correct the conveying error caused by the eccentricity (correction value for eccentricity).\nThe inventors of the present invention, however, have found out that a simple application of the technique disclosed in Japanese Patent Laid-Open No. 2002-273956 has difficulty acquiring a more precise correction value for correcting the conveying error caused by the outer diameter of the conveying roller. If a test pattern with a length equal to an integer multiple of the circumferential length of the roller in the conveying direction is used, a precise correction value of the conveying error can be acquired with the acquiring of correction values for eccentricity and outer diameter being in reverse sequence. In practice, however, it is difficult to form a test pattern with a length that is precisely equal to an integer multiple of the circumferential length of the roller. There has to be, in the test pattern, an excess area that exceeds the area corresponding to an integer multiple of the circumferential length of the roller. Although the correction value for outer diameter can be calculated from the average value of the conveying errors, the part of the above-mentioned excess area must have an influence on the calculation."}
{"text": "Various applications involve sensing (for example measuring) the filtered output current of a switching amplifier, for example the output current provided by an amplifier to a load such as a resistive load through a filter such as an external low-pass for example LC filter.\nA Class D amplifier may be exemplary of such a switching amplifier.\nDespite the extensive activity in that area, improved solutions are desirable in respect of various aspects such as, for example:                avoiding the possible presence of external components such as an external sense resistance with the associated efficiency loss,        effectively addressing technological issues as possibly associated to a load terminal reaching levels above the bridge supply voltage of the amplifier or below ground level,        the possibility of supplying a sense circuit and the output bridges of the amplifier with a same voltage,        possibly filtering the ripple in the bridge current, for example by sensing its mean value, that is the load current, and        achieving high accuracy and precision, for example with offset cancellation.        \nThere is a need in the art to contribute in providing such an improved solution."}
{"text": "1. Field of the Invention\nThe present invention relates to a heat source detector which detects a position of an infrared heat source such as a human body and a moving direction of a heat source position.\n2. Description of the Related Art\nConventionally, radiation (infrared rays) emitted from an infrared heat source such as a human body is detected by an infrared sensor. An infrared sensor detects a heat source by using a passive method in which no light is irradiated onto an object such as a human body to be detected and radiation spontaneously emitted from a heat source is detected. Therefore, the heat source detection can be conducted without adversely affecting the measuring object such as a human body, and without illumination even during the night. Consequently, such a detection using an infrared sensor has a wide variety of applications.\nFIG. 1 shows an example of a conventional heat source detector using an infrared sensor. In the figure, the heat source detector comprises a pyroelectric infrared sensor 1. A shield plate 18 having an aperture 12 is slidably disposed under the infrared sensor 1, and a plurality of Fresnel lenses 16 are arranged under the shield plate 18. In the thus configured detector, an image of radiation from a heat source in each of zones A to C shown in the figure is formed by the Fresnel lenses 16 on an infrared ray sensing device of the infrared sensor 1 so that the heat source is detected. When the aperture 12 of the shield plate 18 is positioned above the zone A as shown in the figure, the zone A constitutes the detection object zone. If a person is in the zone A, for example, an image of radiation from the body of the person is formed on the infrared ray sensing device of the infrared sensor 1, and it is detected that the human body functioning as a heat source is in the zone A. If no person is in the zone A, no radiation is detected. As a result, it is possible to know the presence or absence of a heat source in the zone A.\nWhen the shield plate 18 is slidingly moved so as to be positioned above the zone B and the same operation as that described above is conducted, it is possible in the same manner to know the presence or absence of a heat source in the zone B. Similarly, it is possible also to know the presence or absence of a heat source in the zone C.\nWhen the shield plate 18 is moved at regular time intervals and the presence or absence of a heat source is detected in each of the zones A to C, information can be obtained about the moving direction of a heat source, i.e., whether the heat source is moved from the zone A to the zone C or from the zone C to the zone A.\nAnother example of a conventional heat source detector using an infrared sensor is shown in FIG. 2(a). The detector has a structure wherein a pyroelectric infrared sensor 1 in which eight infrared ray sensing elements 7 are vertically arranged in a row is mounted so as to be reciprocally rotatable as indicated by an arrow A in the figure, a chopper 11 having an aperture 12 is rotatably mounted in the outer peripheral side of the infrared sensor 1, and an infrared lens 10 is disposed in front of the chopper 11. In the detector, radiation from a heat source passes through the aperture 12 of the chopper 11 and its image is formed on the infrared ray sensing elements 7 by the infrared lens 10. The infrared sensor 1 and the chopper 11 are rotated as indicated by the arrow A and an arrow B, respectively, so that radiation is detected while the chopper 11 conducts an opening operation 64 times, whereby the detection of radiation is conducted while changing the detection angle. Resulting signals are processed by a signal processor 24, and a CPU (microcomputer) 25 operates so as to display infrared images of 64.times.8 as shown in FIG. 2(b), thereby detecting the heat source and the moving direction of the heat source.\nToday, visible light CCD cameras are also used in homes, etc. Such a camera is so configured that sensors for detecting visible light are arranged in a high density on a silicon wafer, and high-density wirings corresponding to the sensing devices are formed on the silicon wafer, thereby realizing a high-resolution camera. It is contemplated that the application of this technique to a heat source detector (infrared camera) may enhance the resolution of the heat source detector.\nHowever, a pyroelectric infrared sensor used in a heat source detector is made of ceramics, etc., and hence it is difficult to integrate such sensors on silicon, thereby producing a problem in that, unlike the above-mentioned visible light CCD camera, it is impossible to arrange sensors in a high density on a silicon wafer and fix them to the wafer. It may be contemplated that infrared sensors are arranged by means of bonding so as to arrange them in a high density. However, this involves drawbacks that the bonding operation requires difficult steps, and that steps such as connecting adjacent infrared sensors to each other and drawing out wirings such as lead wires from infrared ray sensing devices of the infrared sensors must be conducted, thereby making the production process very difficult.\nConsequently, a heat source detector such as those shown in FIGS. 1, 2(a) and 2(b) employs a structure wherein infrared sensors are arranged in a low density, and hence its resolution is not very high with the result that such a detector can detect only an approximate position of a heat source. Furthermore, the detector cannot finely determine the moving direction of a heat source. In other words, the detector can conduct only a very rough judgment on the moving direction, i.e., select one of the upper, lower, left, and right directions in the figure, as the moving direction of the heat source.\nIn order to solve the above-mentioned problems, for example, an infrared camera employing quantum infrared sensors which are made of a semiconductor such as HgCdTe or InSb has been proposed, and practically used in some fields. When such a quantum infrared sensor is to be operated, however, the elements of the sensor must be cooled to a very low temperature by using liquid nitrogen or the like, thereby producing problems in that the detector is bulky, and that the production cost of the detector is high."}
{"text": "The present invention relates to a method and apparatus for detecting foreign matters that might be present in liquids. More particularly, the present invention relates to a method and apparatus for detecting undesirable small foreign matters that might be present in medical fluids filled in transparent containers such as injection ampoules and vials, which often contain foreign matters like glass chips, small particles, and fibers. The presence of such foreign matters is not desirable for quality, and it is necessary to select defective containers containing foreign matters by testing all the containers filled with medical fluids.\nIn a conventional method the detection of foreign matters is accomplished as follows: An ampoule to be tested is turned at a high speed and then brought to a standstill quickly. The suspended foreign matters that swirl together with the liquid in the ampoule are illuminated. The beam of light which has passed through the liquid is received by a light detector, and a decrease in light received is regarded as an indication of the presence of foreign matters. In such a method it was difficult to detect foreign matters sensitively with a single photoelectric element as a light detector, because foreign matters to be detected are extremely small as compared with the detection visual field and the difference of photocurrent caused by the presence of foreign matters is also extremely small. Another conventional method in which is used a light detector consisting of many small light sensitive elements equivalent in size to foreign matters and outputs of respective light sensitive elements are scanned, has a fatal drawback that complex electric circuits are required for signal processing and the result of detection is affected by the shape of foreign matters.\nIn the case of the small light sensitive elements having a light sensitive area equal to or smaller than the projected area of the minimum size of foreign matter particles, sensitive detection can be accomplished because the light sensitive surface is shaded by a foreign matter particle and an extremely great difference occurs in quantity of light received between the shaded elements and the unshaded elements. For instance, a light sensitive area of 100.mu..times.100.mu. square will be completely shaded by a particle of about 100.mu..times.100.mu. size, and sensitive detection will be accomplished. However, a long and narrow foreign matter, say 50.mu..times.200.mu. in size, will not cover the square light sensitive surface completely, although the projected area is the same. Therefore, such a foreign matter may not be detected.\nAnother conventional method that detects diffused reflection from foreign matters cannot discriminate foreign matters of different sizes because the ratio of reflection varies depending on the kinds of foreign matters.\nIn the present invention which has been made to overcome the above-mentioned drawbacks, the light sensitive surface of the light detector is divided into a multiplicity of small sections measuring 0.01 mm.sup.2 to 1 mm.sup.2 so that each light sensitive element generates an output signal proportional in magnitude to the projected area of foreign matters, and detection is accomplished by comparing the output signal with the reference value.\nForeign matters vary in shape, and typical shapes are particle and fiber. Foreign matters of fibrous shape generally measure 20.mu. in diameter and more than ten times the diameter in length. Thus, a particle measuring 100.mu..times.100.mu. has the same projected area as a fiber measuring 20.mu..times.500.mu.. If the limit of detection is to be set up for particles measuring 100.mu..times.100.mu. and fibers measuring 20.mu. in diameter (or 500.mu. in length), the light sensitive surface of the detector should be divided into sections measuring 500.mu..times.500.mu. (0.25 mm.sup.2). By measuring the quantity of received light with each section of the divided light sensitive surface, it is possible to obtain an output signal proportional to the projected area of foreign matter regardless of its shape--particle, fiber, and others, and to obtain output signals having a sufficient S/N ratio, without mistaken detection due to small particles inherent to medical fluids.\nAccording to a preferred embodiment of this invention, the light sensitive surface of the light detector is provided with a bundle of optical fibers connected to photoelectric elements in such a manner that each divided section of the light sensitive surface corresponds to each photoelectric element. Thus, the individual sections of the light sensitive surface are substantially continuous and there is no dead zone which might result in failure of detection.\nIn addition, photoelectric converting elements such as phototransistors, photodiodes, and photocells may be staggered directly on the light sensitive surface without using optical fibers, so that foreign matters suspended in the swirling liquid are detected by either row of the detecting elements. Such arrangement eliminates any dead zone.\nIn continuous and automatic detection the visual field of detection should be changed according to the size of ampoules to be tested, and this is accomplished by replacing the light receiver or by covering optically or mechanically a part of the light receiver. Such operation, however, needs skill, and a simple and certain method has been searched for.\nAccording to the present invention the visual field for detection is changed as follows: The output signals from the group of small light sensitive elements are compared with the preset reference values to produce an output. Thus obtained output is then introduced into the detection visual field selector circuit that controls the number of small light sensitive elements to be used according to the size of objects to be detected.\nThe detection visual field selector circuit is a circuit to select a proper number of small light sensitive elements to be used according to the change of visual field relative to the size of objects to be detected. More particularly, it is so designed as to make presetting by means of a selector circuit according to the size of objects to be detected so that the number of small light sensitive elements to be used corresponds to the size of detection visual field. This is accomplished by the built-in matrix circuit such as diode matrix circuit and wired OR-circuit. Thus, it is possible to obtain by simple operation necessary and sufficient outputs from the output signals generated by the small light sensitive elements through the comparators. And the visual field thus changed is extremely accurate."}
{"text": "The general field of the invention is that of data processing, and, more specifically, print stream processing. In its most specific segmentation, the field is that of optimization of those devices directed to processing a print stream for the purpose of producing a plurality of mailpieces.\nIn the past several years, significant changes have occurred in the operation of high volume document production centers. These centers have merged traditional printing capabilities with mailroom production facilities. Executives tasked with the management of both print and mail operations are expected to play an ever-growing role in the creation and design of document centers that will deliver effective, high quality, and high integrity output. The current development and emphasis on these centers in corporations or regional centers has lead to the growing use of the term xe2x80x9cAutomated Document Factoryxe2x80x9d (hereinafter xe2x80x9cADFxe2x80x9d) to describe consolidated printing and mail finishing operations.\nIn current practice, large mailing companies tend to separate the process of creating documents from the process of manufacturing documents (mailpieces). The print center tasked with finishing the created document receives both scheduled and scheduled print jobs with a wide range of requirements. These print jobs are evaluated, scheduled, and executed in the print/finish center.\nBecause the print/finish center has traditionally been xe2x80x9cinformation systems poor,xe2x80x9d most of the work required to prepare or xe2x80x9cconditionxe2x80x9d the print job for manufacturing was created in the business unit or print service client. Typical conditioning processes include: performing postal address hygiene; adding PostNet(trademark) barcodes; presorting mailings; adding inserter barcode instructions; adjusting printer paper size and orientation; and, adding spot color instructions.\nThe manager of such a print/finish operation, seeking to maximize efficiency through optimal use of equipment and decision making tools, is faced with a dilemma. First, the decisions about the structure and management of the print/finish center are generally made outside of the center; the decisions are generally made by the Information Systems (IS) group creating the print job and its associated print stream. Document manufacturing requests are also assigned lower priorities, further limiting management control. Second, the hardware systems and their associated peripheral devices are often sourced from different manufacturers so that the printers and inserters being fed by the print stream are relying on differing motivators from the print stream.\nTo help classify and organize the concept of the emerging print/finish center, an architecture has been developed within the print stream industry that is referred to as the ADF. The Automated Document Factory(trademark) architecture proposed by the Gartner Group of Stamford, Conn., provides a model for a set of processes that prepares and positions enterprises to manage the creation and delivery of high-volume digitized documents by using factory production techniques that appropriately and optimally mechanize document production. The raw materials of production (i.e., the document data and preparation instructions), enter the ADF which transforms them into digital documents and prepares them for delivery.\nThe architecture for the ADF is comprised of four (4) modules; these include: input; transformation; delivery and preparation; and, control and reporting. Each module, or building block, is made up of other modules and each is connected by a series of interfaces, or links.\nEach of the building blocks must be linked through effective communication which includes the tracking and measurement of the input and output of the document manufacturing hardware and associated peripherals. To enhance productivity and costeffectiveness of the overall system, systems managers need to be able to scrutinize every element of the print job process to see where improvements can be made. Thus, each of the modules takes on an increased significance when viewed with respect to their relationship with the overall system.\nThere is thus a need to provide each of the modules for the ADF so that the structure can be self supporting and viable. The input module is where all of the data and instructions needed to transform the arriving print stream data into documents enters the ADF. The present invention is currently being introduced to the print stream market by the assignee of the present invention, Pitney Bowes Inc. of Stamford, Conn., as the InStream(trademark) server which is designed as the input module for the ADF.\nIt is an object of the present invention to provide the input module to the conceptual ADF frame by describing herein an open systems, client-server technology for facilitating automated document manufacturing techniques.\nIt is a further object of the present invention to provide a method of optimizing the use of hardware and associated peripheral devices in manufacturing documents that have been digitally delivered through the input module. Additionally, it is further object of the present invention to measure the activities of each of the hardware and peripheral components so that accurate reporting can be made so as to facilitate subsequent job performance decisions and so as to maximize system utility and performance.\nAccording to the invention, the above objects are achieved and the disadvantages of the prior art are overcome by a method and system for print stream determination wherein the print stream is directed to a print processing application for determining a job to be performed by a print processing system. The preferred embodiment of the system will include peripheral devices for producing mailpieces from the print stream as determined by job tickets established by the print processing application.\nThe system\"\"s method begins with the initiation of a print stream processing application to which a print stream is directed. The print stream is initiated at any one of a number of print service clients which would include: print servers; digital document delivery; and, finishing equipment. A print job is determined from a set of characteristics resident in the print stream; these characteristics might include: document text; mailpiece definition; volume; and, routing instructions for routing the print stream to a corresponding peripheral device driver of a corresponding peripheral device; and, elements for determining device driver definition elements for defining characteristics and capabilities of each of the corresponding peripheral devices.\nThe optimal use determination is made by comparing each of the job definition elements with each of the device driver definition elements. The comparison of the job definition elements with each of the device driver definition elements further comprises a number of steps which begins with determination of a range or a value for each of the job definition elements wherein the range or the value is representative of a desired result. A range or a value for each of the device driver definition elements is determined wherein the range or value is representative of a potential result. Each of the ranges or values for the device driver definition elements is compared to each corresponding one of the ranges or values for the job definition elements to determine whether or not each of the desired results has a corresponding potential result.\nIf no corresponding potential result can be established in the comparison steps of the methodology, then the application queries the print processing system to determine whether or not an alternative peripheral device, with a matching value or range, is a component of the system. If an alternative peripheral device is a component of the system, then the application performs a comparison of the non-corresponding potential result with the matching value for the alternative peripheral; and, if no peripheral device is a component of the system, then the application will select a default peripheral device. The print stream will then be directed to the selected device for processing. The optimal peripheral devices might include: a printer; a second memory; a data processing application for further applying the application output stream; an envelope printer for printing address data on one or more envelopes; a document printer for printing one or more pages; a collator for collating a plurality of pages printed by the document printer; and/or, an inserter for inserting the one or more printed pages into a corresponding envelope.\nBased upon print stream characteristics, the print processing application will determine the optimal use of the system\"\"s peripheral devices for performing the job. Additionally, the application will determine one or more application output streams wherein each one of the one or more output streams is directed to a corresponding peripheral device. Once the print stream has been analyzed, and the use of peripheral devices determined, the application establishes a job ticket representative of the print job and the optimal use of the one or more corresponding peripheral devices. The job ticket is saved in a memory of the print processing system for current and/or later use. If the job is to be performed essentially simultaneously with the processing of the stream, then the system performs the print job as described in saved job ticket; and can produce one or more reports indicative of the job performance. The reports are printed on a printer operatively associated with the print processing system, or may be downloaded to a memory medium such as a diskette, tape, or hard disk, or to a communication link such as a modem, or to both a memory and a communications device.\nThe method is capable of performing other tasks as derived from a legacy application which contributes a set of print data to the print processing system. The print processing system applies the set of print data to a set of one or more pre-determined job processes which might include: postal sorting; address cleansing; and, bar coding."}
{"text": "1. Field of the Invention\nThe present invention relates to an optical head apparatus and an optical information recording and reproducing apparatus that record and reproduce information to and from an optical recording medium, in particular, to those that allow the deviation of the thickness of the substrate of an optical recording medium and the radial tilt thereof to be detected without offsets of a focusing error signal and a tracking error signal.\n2. Description of the Related Art\nOn a write-once type optical recording medium and a rewritable optical recording medium on which an RF (Radio Frequency) signal has not been recorded, grooves have been formed as tracks. Generally, viewed from the incident light side of such an optical recording medium, a concave portion and a convex portion of a groove are referred to as land portion and groove portion, respectively. When a focusing error signal of a write-once type optical recording medium or a rewritable optical recording medium is detected, the signal level of the focusing error signal does not strictly accord with the defocus amount. In other words, at the position where the defocus amount is zero, the signal level of the focusing error signal is not strictly zero. Since grooves have been formed, theoretically, the sign of an offset of the focusing error signal detected at the land portion is reverse of that detected at the groove portion. The offset is referred to as groove traverse noise offset. When a tracking error signal is detected from a write-once type optical recording medium or a rewritable optical recording medium, push-pull method is normally used. However, in the push-pull method, when an objective lens of an optical head apparatus is shifted in the radial direction of an optical recording medium, the tracking error signal has an offset. This offset is referred to as lens shift offset. To prevent recording and reproduction characteristics from deteriorating due to such offsets, it is necessary for an optical head apparatus and an optical information recording and reproducing apparatus to provide a mechanism that allows the focusing error signal and the tracking error signal to have no offsets.\nFIG. 1 shows the structure of a conventional optical head apparatus that allows the focusing error signal and the tracking error signal to have no offsets. The optical head apparatus has been disclosed as Japanese Patent Laid-Open Publication No. 2000-82226. Light emitted from a semiconductor laser 1 is separated into three pieces of 0-th order light and ±1st order diffracted light by a diffractive optical element 3p. Around 50% of these three pieces of light transmits a beam splitter 24. A collimator lens 2 collimates those pieces of light. An objective lens 6 focuses those pieces of light on a disk 7. Three pieces of light that are reflected from the disk 7 transmit the objective lens 6 and the collimator lens 2 in the reverse direction. The beam splitter 24 reflects around 50% of those piece of light. Those pieces of light transmit a cylindrical lens 8. A photo detector 10b receives those pieces of light. The photo detector 10b is disposed at the middle position of two focal lines of the collimator lens 2 and the cylindrical lens 8.\nFIG. 2 shows the positions of focused spots on the disk 7. Focused spots 13a, 13t, and 13u correspond to 0-th order light, +1st order diffracted light, and −1st order diffracted light that are separated by the diffractive optical element 3p. The focused spot 13a is placed on a track 12 (land portion or groove portion). The focused spot 13t is placed on a right adjacent track (groove portion or land portion) of the track 12. The focused spot 13u is placed on a left adjacent track (groove portion or land portion) of the track 12.\nFIG. 3 shows the pattern of a light receiving portion of the photo detector 10b and the positions of light spots on the photo detector 10b. A light spot 15a corresponds to 0-th order light separated by the diffractive optical element 3p. The 0-th order light is received by four light receiving portions 25a to 25d divided by a division line that passes through the optical axis and that parallels the tangential direction of the disk 7 and a division line that parallels the radial direction thereof. A light spot 15j corresponds to +1st order diffracted light separated by the diffractive optical element 3p. The light spot 15j is received by four light receiving portions 25e to 25h divided by a division line that passes through the optical axis and that parallels the tangential direction of the disk 7 and a division line that parallels the radial direction thereof. A light spot 15k corresponds to −1st order diffracted light separated by the diffractive optical element 3p. The light spot 15k is received by four light receiving portions 25i to 25l divided by a division line that passes through the optical axis and that parallels the tangential direction of the disk 7 and a division line that parallels the radial direction thereof. The focused spots 13a, 13t, and 13u are trained nearly in the tangential direction of the disk 7. In contrast, the light spots 15a, 15j, and 15k are trained nearly in the radial direction due to the operations of the collimator lens 2 and the cylindrical lens 8.\nWhen outputs of the light receiving portions 25a to 25l are denoted by V25a to V25l, respectively, the focusing error signal can be obtained by differential astigmatism method as (V25a+V25d)−(V25b+V25c)+K{(V25e+V25h+v25i+V25l)−(V25f+V25g+V25j+V25k)} (where K is constant). On the other hand, the tracking error signal can be obtained by differential push-pull method as (V25a+V25b)−(V25c+V25d)−K{(V25e+V25f+V25i+V25j)−(V25g+V25h+V25k+V25l)}. In addition, an RF signal of the focused spot 13a can be obtained as V25a+V25b+V25c+V25d . \nFIGS. 4A, 4B, and 4C show various types of focusing error signals. In FIGS. 4A to 4C, the horizontal axis represents the defocus amount of the disk 7, whereas the vertical axis represents each focusing error signal. A focusing error signal 26a shown in FIG. 4A is a focusing error signal of the focused spot 13a that is placed on a land portion. A focusing error signal 26b shown in FIG. 4A is a focusing error signal of the focused spot 13a placed on a groove portion. A focusing error signal 26c shown in FIG. 4B is a focusing error signal of the focused spots 13t and 13u in the case that the focused spot 13a is placed on a land portion. A focusing error signal 26d shown in FIG. 4B is a focusing error signal of the focused spots 13t and 13u in the case that the focused spot 13a is placed on a groove portion. The signal level of a focusing error signal is not strictly zero at the position where the defocus amount is zero. In FIG. 4A, a focusing error signal has a positive offset at a land portion, whereas it has a negative offset at a groove portion. In FIG. 4B, a focusing error signal has a negative offset at a land portion, whereas it has a positive offset at a groove portion. In contrast, a focusing error signal 26e shown in FIG. 4C is a final focusing error signal that is the sum of the focusing error signal of the focused spot 13a and the focusing error signals of the focused spots 13t and 13u in the case that the focused spot 13a is placed at a land portion and a groove portion. In FIG. 4C, the offsets of the focusing error signals shown in FIGS. 4A and 4B are cancelled each other. Thus, the focusing error signal 26e does not have an offset.\nFIGS. 5A, 5B, and 5C show various types of tracking error signals. In FIGS. 5A, 5B, and 5C, the horizontal axis represents the off track amount of the disk 7, whereas the vertical axis represents each tracking error signal. A tracking error signal 27a shown in FIG. 5A is a tracking error signal of the focused spot 13a in the case that the objective lens 6 is outward shifted in the radial direction of the disk 7. A tracking error signal 27b shown in FIG. 5A is a tracking error signal of each of the focused spots 13t and 13u in the case that the objective lens 6 is outward shifted in the radial direction of the disk 7. A tracking error signal 27c shown in FIG. 5B is a tracking error signal of the focused spot 13a in the case that the objective lens 6 is inward shifted in the radial direction of the disk 7. A tracking error signal 27d shown in FIG. 5B is a tracking error signal of each of the focused spots 13t and 13u in the case that the objective lens 6 is inward shifted in the radial direction of the disk 7. The polarity of the tracking error signal of the focused spot 13a is the reverse of the polarity of the tracking error signal of each of the focused spots 13t and 13u. However, when the objective lens 6 is shifted in the radial direction of the disk 7, the sign of the offset of the tracking error signal of the focused spot 13a is the same as the sign of the offset of the tracking error signal of each of the focused spots 13t and 13u. In FIG. 5A, the offsets of the tracking error signal 27a and the tracking error signal 27b are positive. In FIG. 5B, the offsets of the tracking error signals 27c and 27d are negative. In contrast, a tracking error signal 27e shown in FIG. 5C is a final tracking error signal that is the difference between a tracking error signal of the focused spot 13a and a tracking error signal of each of the focused spots 13t and 13u in the case that the objective lens 6 is outward and inward shifted in the radial directions of the disk 7. In FIG. 5C, the offsets of the tracking error signals shown in FIGS. 5A and 5B are cancelled each other. Thus, the tracking error signal 27e shown in FIG. 5C does not have an offset.\nThe recording density of an optical information recording and reproducing apparatus is reversely proportional to the second power of the diameter of a focused spot formed on a record medium by an optical head apparatus. In other words, as the diameter of a focused spot is smaller, the recording density becomes higher. The diameter of a focused spot is reversely proportional to the numerical aperture of an objective lens of the optical head apparatus. In other words, as the numerical aperture of the objective lens is higher, the diameter of a focused spot becomes smaller. On the other hand, when the thickness of the substrate of the optical recording medium deviates from a designed value, the shape of a focused spot deforms due to the spherical aberration caused by the deviation of the thickness of the substrate. As a result, the recording and reproduction characteristics deteriorate. Since the spherical aberration is proportional to the fourth power of the numerical aperture, as the numerical aperture of the objective lens is higher, the margin of the deviation of the thickness of the substrate of the optical recording medium becomes narrower. When the optical recording medium tilts in the radial direction against the objective lens, the shape of a focused spot deforms due to the comatic aberration caused by the radial tilt. As a result, the recording and reproduction characteristics deteriorate. Since the comatic aberration is proportional to the third power of the numerical aperture of the objective lens, as the numerical aperture of the objective lens is higher, the margin of the radial tilt of the optical recording medium against the recording and reproduction characteristics becomes narrower. Thus, in the optical head apparatus and the optical information recording and reproducing apparatus that use an objective lens having a higher numerical aperture for a higher recording density, it is necessary to detect and compensate the deviation of the thickness of the substrate of the optical recording medium and the radial tilt thereof so as to prevent the recording and reproduction characteristics from deteriorating.\nAn example of a conventional optical head apparatus that can detect the radial tilt of an optical recording medium has been described in SPIE Proceedings, Vol. 4090, pp. 309–318. The optical head apparatus features a diffractive optical element 3q with which the diffractive optical element 3p of the optical head apparatus shown in FIG. 1 is substituted.\nFIG. 6 is a plan view showing the diffractive optical element 3q. The diffractive optical element 3q has a structure in which a diffraction grating is formed in only the inside of a circular area 28 of which diameter is smaller than the effective diameter (denoted by a dotted line in FIG. 6) of the objective lens 6. The direction of grating members of the diffraction grating parallels the radial direction of the disk 7. The grating members are linearly formed and equally spaced. Light that enters the inside of the circular area 28 partly transmits it as 0-th order light and is diffracted as +1st order diffracted light. In contrast, light that enters the outside of the circular area 28 fully transmits it. In other words, 0-th order light separated by the diffractive optical element 3q contains both light that transmits the inside of the circular area 28 and light that transmits the outside of the circular area 28. Thus, the numerical aperture of 0-th order light depends on the effective diameter of the objective lens 6. On the other hand, +1st order diffracted light separated by the diffractive optical element 3q contains only light diffracted on the inside of the circular area 28. Thus, the numerical aperture of ±1st order diffracted light depends on the diameter of the circular area 28. As a result, the distribution of the intensity of 0-th order light separated by the diffractive optical element 3q is different from that of each of +1st order diffracted light separated thereby. Thus, the intensity of the peripheral portion of each of +1st order diffracted light is lower than that of 0-th order light.\nFIG. 7 shows the positions of focused spots on the disk 7. Focused spots 13a, 13v, and 13w correspond to 0-th order light, +1st order diffracted light, and −1st order diffracted light separated by the diffractive optical element 3q, respectively. The focused spot 13a, 13v, and 13w are placed on the same track 12 (land portion or groove portion). Since the intensity of the peripheral portion of each of ±1st order diffracted light is lower than that of 0-th order light, the diameter of each of the focused spots 13v and 13w that are +1st order diffracted light is larger than the diameter of the focused spot 13a that is 0-th order light.\nThe pattern of the light receiving portion of the photo detector of the optical head apparatus and the positions of the light spots on the photo detector are the same as those shown in FIG. 3.\nFIGS. 8A, 8B, and 8C show various types of tracking error signals for detecting a radial tilt. In FIGS. 8A, 8B, and 8C, the horizontal axis represents the off track amount of the disk 7, whereas the vertical axis represents each tracking error signal. A tracking error signal 29a shown in FIG. 8A is a tracking error signal of each of the focused spots 13a, 13v, and 13w in the case that the disk 7 does not have a radial tilt. In contrast, a tracking error signal 29b shown in FIG. 8B is a tracking error signal of the focused spot 13a in the case that the disk 7 has a positive radial tilt. A tracking error signal 29c shown in FIG. 8B is a tracking error signal of each of the focused spots 13v and 13w in the case that the disk 7 has a positive radial tilt. A tracking error signal 29d shown in FIG. 8C is a tracking error signal of the focused spot 13a in the case that the disk 7 has a negative radial tilt. A tracking error signal 29e shown in FIG. 8C is a tracking error signal of each of the focused spots 13v and 13w in the case that the disk 7 has a negative radial tilt. The position at which the tracking error signal of the focused spot 13a traverses zero from the − side to the + side as the curve on the graph moves rightward corresponds to a land portion, whereas the position at which the tracking error signal of the focused spot 13a traverses zero from the + side to the − side as the curve on the graph moves rightward corresponds to a groove portion. When the disk 7 does not have a radial tilt, the zero crossing point of the tracking error signal of each of the focused spots 13v and 13w matches the zero crossing point of the tracking error signal of the focused spot 13a. Thus, the signal level of each of the tracking error signals becomes zero regardless of whether each of the focused spots 13v and 13w is placed on a land portion or a groove portion. When the disk 7 has a positive radial tilt, the zero crossing point of the tracking error signal of each of the focused spots 13v and 13w shifts leftward against the tracking error signal of the focused spot 13a. Thus, when each of the focused spots 13v and 13w is placed at a land portion, the signal level of the tracking error signal becomes positive. In contrast, when each of the focused spots 13v and 13w is placed at a groove portion, the signal level of the tracking error signal becomes negative. When the disk 7 has a negative radial tilt, the zero crossing point of the tracking error signal of each of the focused spots 13v and 13w shifts rightward against the tracking error signal of the focused spot 13a. Thus, when each of the focused spots 13v and 13w is placed at a land portion, the signal level of the tracking error signal becomes negative. In contrast, when each of the focused spots 13v and 13w is placed at a groove portion, the signal level of the tracking error signal becomes positive. Thus, the tracking error signal of each of the focused spots 13v and 13w when performing a tracking servo using the tracking error signal of the focused spot 13a can be used as a radial tilt signal.\nIn the conventional optical head apparatus that allows the focusing error signal and the tracking error signal to have no offsets, the sum of a focusing error signal of 0-th order light separated by a diffractive optical element and a focusing error signal of each of ±1st order diffracted light separated by the diffractive optical element is a final focusing error signal. In addition, the difference between a tracking error signal of 0-th order light and a tracking error signal of each of ±1st order diffracted light is a final tracking error signal. To cancel an offset of a focusing error signal of 0-th order light separated by a diffractive optical element and that of each of ±1st order diffracted light separated thereby and to cancel an offset of a tracking error signal of 0-th order light separated by the diffractive optical element and that of each of ±1st order diffracted light separated thereby, it is necessary to cause the distribution of the intensity of 0-th order light separated by the diffractive optical element to be the same as that of each of ±1st order diffracted light separated thereby.\nOn the other hand, in the conventional optical head apparatus that can detect a radial tilt of an optical recording medium, a radial tilt thereof is detected corresponding to the deviation between the zero crossing point of a tracking error signal of 0-th order light separated by the diffractive optical element and the zero crossing point of the tracking error signal of each of ±1st order diffracted light separated thereby, so it is necessary to cause the distribution of the intensity of 0-th order light separated by the diffractive optical element to be different from that of each of ±1st order diffracted light separated thereby.\nThus, in the conventional optical head apparatus, both the structure that allows a focusing error signal and a tracking error signal to have no offsets and the structure that allows a radial tilt of an optical recording medium to be detected cannot be satisfied at a time."}
{"text": "1. Field of the Invention\nThe invention relates to a charge injection device (CID) for sensing infrared (IR), and more particularly to an improved readout circuit having plural readout channels optimized for a cryogenically operated IR sensor.\n2. Prior Art\nInfrared sensors, initially having smaller numbers of pixel sites, are now being developed with larger numbers of sites, comparable to visible light sensors. Two dimensional arrays of 128.times.128 and 256.times.256 sites are in development, as well as linear arrays having as many as 128.thrfore.12, or 128.times.24 sites. Both large two dimensional IR arrays and large linear IR arrays require parallel readout circuitry. Parallel readout circuitry is required both for greater speed in sampling the sites and for greater optimization of array sensitivity.\nThe problem presented in providing parallel readout circuitry is that of not interfering with the cryogenic environment, which is required for high performance infrared focal plane arrays, while at the same time avoiding crosstalk between parallelled channels. The focal plane array customarily combines an IR detector chip and readout circuitry. The readout circuitry is required to convert the signal obtained by sampling individual IR detector elements on the chip into a (usually) single serial signal. The signal must be amplified at each stage of the processing to a sufficiently high level to avoid noise problems in subsequent processing.\nInterference with the cryogenic performance may occur if large numbers of high conductance electrical conductors leading from the cryogenic environment of the focal plane array to the warmer ambient environment are required for parallel readout. Any additional electrically conductive paths will provide additional thermally conductive paths which will add to the thermal load presented when the focal plane array is cooled. Added thermal bulk is undesirable for similar reasons.\nThe conventional solution to crosstalk between parallelled channels tends to compromise cryogenic efficiency. In particular, the requirements for thermal isolation make it difficult to provide Dewar penetrations of low electrical impedance that would reduce crosstalk. In addition, focal plane miniaturization prohibits adding the bulk required for effective electrical filtering on the focal plane that would reduce crosstalk. Consequently, crosstalk of the focal plane signals between channels due to voltage modulation in the power supply paths tend to be present if cryogenic efficiency is optimized."}
{"text": "Manufacture of wood composite board and other building materials requires the evaluation of many variables. The size of the wood fibers and particles, the resin or adhesive content, the pressure, the pressing time, the temperature, and other factors must be determined and checked in order to produce a board of consistent quality.\nAn accepted measure of consistency is the density profile of the board. If the density profile of a composite board with the desired strength characteristics can be reproduced, then the strength characteristics will be maintained.\nSeveral different methods for evaluating density profiles have been used in industry and in the laboratory. The gravimetric method, in which the specimen is planed, measured, and weighed in small increments, is the most direct. The density measurements, in this case, can be evaluated approximately every 0.030 inches (0.076 cm) of thickness.\nTwo methods involving radiation sources have also been evaluted. One such method uses X rays to produce a photographic image of the profile of a composite wood specimen; this method allows for virtually continuous density measurements across the cross-section of the board, with correlation made with known standards. A second such method uses a gamma ray source to scan the thickness of a specimen. The transmitted radiation is then correlated to the material density, and a continuous measurement is again possible.\nOf the methods mentioned and discussed above, the conventional gravimetric method for density determination is the most time-consuming and the least resolute. Although the density calculations are the most direct (weight/volume), the measurements become more difficult to obtain as the specimen is planed, because the thickness decreases. A thicker specimen would be easier to handle, but would also require more time, due to the greater number of sections that would have to be removed.\nThe X ray and gamma ray methods for density measurement each provide continuous profiles of the specimen, but they also require radiation sources and a greater knowledge of equipment operation by the user. These methods are faster than the gravimetric method and provide information about densities at specific points rather than averaging the densities over small sections. Both such methods claim to be nondestructive: However, this would not be the case in practical applications in industry. The samples being evaluated must be no more than 3 inches (7.62 cm) in width; so some part of the manufactured board must be cut. The operation of these densitometers not only requires training: The densitometers also incur a substantial initial cost before consistent measurements can be made."}
{"text": "Currently available automated liquid soap dispensers can deliver liquid soap automatically in response to the presence of a nearby object, such as a user's hand. These liquid soap dispensers typically employ an infrared sensor to detect the nearby object. Upon detection of the nearby object, the infrared sensor sends a signal to activate a motor, which in turn drives a shaft which drives a pump. The pump pumps the liquid soap to a spout of the dispenser which dispenses the soap. Typically, the motor is located outside of the dispenser, whereas, the pump is submerged in the liquid soap within the dispenser. The drive shaft penetrates the dispenser through an orifice typically at the base of the dispenser. A seal is typically provided surrounding the shaft and sealing the orifice. However, with use the seal wears out and the liquid soap leaks out of the dispenser. Often, the liquid soap leaks onto the motor or the motor circuitry causing failure of the motor.\nFurthermore, most liquid dispensers have reservoirs which are mounted below a countertop. Consequently, accessing of the reservoir for refilling with liquid soap is inconvenient. Typically a reservoir of the soap dispenser needs to be removed from below the counter so that it may be filled. When removed, tubing which is used to deliver the liquid soap to the spout is exposed and liquid soap on such tubing drips on the surrounding surfaces. Moreover, with some soap dispenser, the motor may have to be removed before the reservoir is removed for refilling. Thus, a soap dispenser is desired that overcomes the aforementioned problems."}
{"text": "1. Field of the Invention\nThe present invention relates to an apparatus and method using expandable tubulars to complete a well. More particularly, the invention relates to the installation of an expandable sand screen. More particularly still, the invention relates to a single trip installation process to set a liner hanger in a wellbore and then expand a sand screen.\n2. Description of the Related Art\nHydrocarbon wells are typically formed with a central wellbore that is supported by steel casing. The casing lines a borehole formed in the earth during the drilling process. An annular area formed between the casing and the borehole is filled with cement to further support and form the wellbore.\nSome wells are produced by perforating the casing of the wellbore at selected depths where hydrocarbons are found. Hydrocarbons migrate from the formation through the perforations and into the wellbore where they are usually collected in a separate string of production tubing for transportation to the surface of the well. In other instances, a lower portion of a wellbore is left open and not lined with casing. This xe2x80x9copen holexe2x80x9d completion permits hydrocarbons in an adjacent formation to migrate directly into the wellbore where they are subsequently raised to the surface, possibly through an artificial lift system.\nOpen hole completions can provide higher production than cased hole completions and they are frequently utilized in connection with horizontally drilled boreholes. However, open hole completions leave aggregate material, including sand, free to invade the wellbore. Sand entering an open hole wellbore causes instability within the open hole which enhances the risk of complete collapse. Sand production can also result in premature failure of artificial lift and other downhole and surface equipment due to the abrasive nature of sand. In some instances, high velocity sand particles can contact and erode lining and tubing.\nSand can also be a problem where casing is perforated to collect hydrocarbons. Typically, casing is perforated with a perforating assembly or xe2x80x9cgunsxe2x80x9d that are run into a wellbore and fired to form the perforations. Thereafter, the assembly is removed and a separate assembly is installed to collect the migrating hydrocarbons. The perforations also create a passageway for aggregate material, including sand to enter the wellbore. As with an open wellbore, sand entering the cased wellbore can interfere with the operation of downhole tools, clog screens and damage components, especially if the material enters the wellbore at a high velocity.\nTo control particle flow into a wellbore, well screens are often employed downhole. Conventional wellscreens are placed adjacent perforations or unlined portions of the wellbore to filter out particulates as production fluid enters a tubing string. One form of well screen recently developed is the expandable sand screen (ESS). In general, the ESS is constructed of different composite layers, including a filter media.\nA more particular description of an ESS is found in U.S. Pat. No. 5,901,789, which is incorporated by reference herein in its entirety. That patent describes an ESS which consists of a perforated base pipe, a woven filtering material, and a protective, perforated outer shroud. Both the base pipe and the outer shroud are expandable, and the woven filter is typically arranged over the base pipe in sheets that partially cover one another and slide across one another as the sand screen is expanded. The ESS is expanded by a cone-shaped object urged along its inner bore or by an expander tool having radially outward extending rollers that are fluid powered from a tubular string. Using expansion means like these, the ESS is subjected to outwardly radial forces that urge the expanding walls against the open formation or parent casing. The components of the ESS are expanded past their elastic limit, thereby increasing the inner and outer diameter of the tubular.\nA major advantage to the ESS in an open wellbore is that once expanded, the walls of the wellbore are supported by the ESS. Additionally, the annular area between the screen and the wellbore is mostly eliminated, and with it the need for a gravel pack. A gravel pack is used with conventional well screens to fill an annular area between the screen and wellbore and to support the walls of the open hole. With an ESS, the screen is expanded to a point where its outer wall places a stress on the walls of the wellbore, thereby providing support to the walls of the wellbore to prevent dislocation of particles. Solid expandable tubulars are oftentimes used in conjunction with an ESS to provide a zonal isolation capability. In addition to open wellbores, the ESS is effectually used with a perforated casing to control the introduction of particulate matter into the cased wellbore via the perforations.\nWhile an ESS can reduce or eliminate the inflow of particles into a wellbore, the screen must be installed and expanded in order to operate effectively. Any delay in the installation permits additional time for sand to enter the wellbore and the time period is especially critical between the formation of perforations in a casing wall and the expansion of screen against the perforations. The delays are especially critical if the newly formed wellbore is placed in an over balanced condition prior to expanding the ESS. An overbalanced condition permits fluids to enter the formations and hamper later production of hydrocarbons.\nIn current installation procedures of ESS the operator makes two trips downhole. In the first trip, the operator sets a liner hanger to secure the ESS in the wellbore. After returning from the first trip downhole, the operator must make a second trip with an expansion tool in order to expand the ESS.\nThere are several disadvantages to a multiple trip installation procedure. The biggest disadvantage relates to expensive downtime necessary to make both trips. Also, a delay between the first and second trips can cause well control problems due to fluid loss. For example, pressurized fluid in the wellbore used to actuate various mechanical components during the installation process can enter the formations causing formations to clog-up or collapse, restricting the flow of hydrocarbons. In addition, loss of drilling fluid increases the completion cost of the well. In other instances, a delay between perforating a casing and expanding a sand screen against the perforations increases the likelihood that solids from the formations will enter the wellbore. In addition to the foregoing, packers used to fix an ESS in a wellbore often have a relatively small inside diameter. These packer-like components remain in the wellbore and can cause access problems for remedial work required below the suspension device.\nThere is a need therefore, for an apparatus to reduce the time needed to install an expandable sand screen in a wellbore. There is a further need to set a sand screen in a wellbore and then expand the sand screen in a single trip. There is a further need for a method and apparatus to facilitate the setting of a liner hanger in a wellbore prior to the expansion of an ESS. Still further, there is a need for an apparatus to minimize the exposure to formation solids before expanding the ESS. There is a further need for a single trip ESS apparatus that uses a liner hanger that does not restrict access within the wellbore after the ESS is expanded.\nThe present invention includes a method and apparatus for installing and expanding an ESS in a wellbore in a single trip. In one aspect of the invention, a liner hanger and expandable screen are provided and are run into the wellbore with an expansion tool and work string. After the hanger is set, the expansion tool is used to expand the screen. In another aspect, an annular area within the apparatus is utilized in order to set the hanger with pressurized fluid. Thereafter, cup packers used in sealing the annulus are lifted from the liner prior to expanding the screen. The expansion tool and work string are then removed leaving the expanded ESS and hanger in the wellbore."}
{"text": "A conventional bandgap voltage reference circuit is based on the addition of two voltage components having opposite and balanced temperature slopes.\nFIG. 1 illustrates a symbolic representation of a conventional bandgap reference. It consists of a current source, 110, a resistor, 120, and a diode, 130. It will be understood that the diode represents the base-emitter junction of a bipolar transistor. The voltage drop across the diode has a negative temperature coefficient, TC, of about −2.2 mV/° C. and is usually denoted as a Complementary to Absolute Temperature (CTAT) voltage, since its output value decreases with increasing temperature. This voltage has a typical negative temperature coefficient according to equation 1 below:\n                                          V            be                    ⁡                      (            T            )                          =                                            V                              G                ⁢                                                                  ⁢                0                                      ⁡                          (                              1                -                                  T                                      T                    0                                                              )                                +                                                    V                be                            ⁡                              (                                  T                  0                                )                                      *                          T                              T                0                                              -                      σ            *                          KT              q                        *            ln            ⁢                                                  ⁢                          (                              T                                  T                  0                                            )                                +                                    KT              q                        *                          ln              (                                                Ic                  ⁡                                      (                    T                    )                                                                    Ic                  ⁡                                      (                                          T                      0                                        )                                                              )                                                          (                  Eq          .                                          ⁢          1                )            Here, VG0 is the extrapolated base-emitter voltage at zero absolute temperature, of the order of 1.2V; T is actual temperature; T0 is a reference temperature, which may be room temperature (i.e. T=300K); Vbe(T0) is the base-emitter voltage at T0, which may be of the order of 0.7V; σ is a constant related to the saturation current temperature exponent, which is process dependent and may be in the range of 3 to 5 for a CMOS process; K is the Boltzmann's constant, q is the electron charge, Ic(T) and Ic(T0) are corresponding collector currents at actual temperatures T and T0, respectively.\nThe current source 110 in FIG. 1 is desirably a Proportional to Absolute\nTemperature (PTAT) source, such that the voltage drop across resistor 120 is PTAT voltage. As absolute temperature increases, the voltage drop across resistor 120 increases as well. The PTAT current is generated by reflecting across a resistor a voltage difference (ΔVbe) of two forward-biased base-emitter junctions of bipolar transistors operating at different current densities. The difference in collector current density may be established from two similar transistors, i.e. Q1 and Q2 (not shown), where Q1 is of unity emitter area and Q2 is n times unity emitter area. The resulting ΔVbe, which has a positive temperature coefficient, is provided in equation 2 below:\n                              Δ          ⁢                                          ⁢                      V            be                          =                                                            V                be                            ⁡                              (                                  Q                  1                                )                                      -                                          V                be                            ⁡                              (                                  Q                  2                                )                                              =                                    KT              q                        *                          ln              ⁡                              (                n                )                                                                        (                  Eq          .                                          ⁢          2                )            \nIn some applications, for example low power applications, the resistor 120 may be large and even dominate the silicon die area, thereby increasing cost. Therefore, it is desirable to have PTAT voltage circuits which are resistorless. PTAT voltages generated using active devices may be sensitive to process variations, via offsets, mismatches, and threshold voltages. Further, active devices used in PTAT voltage cells may contribute to the total noise of the resulting PTAT voltage. One goal of an embodiment of the present invention is to provide a resistorless PTAT cell operable at low power with little sensitivity to process variations and having low noise.\nFIG. 2 illustrates the operation of the circuit of FIG. 1. By combining the CTAT voltage, V_PTAT of diode 130 with the PTAT voltage, V_PTAT, from the voltage drop across resistor 120, it is possible to provide a relatively constant output voltage Vref over a wide temperature range (i.e. −50° C. to 125° C.). This base-emitter voltage difference, at room temperature, may be of the order of 50 mV to 100 mV, for n from 8 to 50.\nTo balance the voltage components of the negative temperature coefficient from equation 1 and the positive temperature coefficient of equation 2, it is desirable to have the capability of fine-tuning the PTAT component to improve the immunity to process variations. Accordingly, in another embodiment of the present invention, a goal is to provide a fine-tune capability of the PTAT component.\nIn yet another embodiment of the present invention, it is a goal to multiply the ΔVbe component of transistors which are operated at different current densities to provide a higher reference voltage which is insensitive to temperature variations."}
{"text": "1. Field of the Invention\nThis invention relates to an Infiniband™ (IB) input/output unit (IOU) that enables a host computing system to access Input/Output (I/O) devices that conform to the Small Computer System Interface (SCSI) mapping onto Fibre Channel (FC). Such I/O devices are referred to as FCP I/O devices.\n2. Description of Background\nThe following references pertain to this invention:                1. Information regarding IB subnet management and administration can be found in Chapters 14, 15 and 16 of Infiniband Architecture Specification, Volume 1, release 1.1, Infiniband Trade Association. Infiniband is a trademark of SYSTEM I/O doing business as InfiniBand Trade Association of Portland Oreg.        2. Information regarding Fibre Channel can be found in Fibre Channel-Framing and Signalling (FC-FS) rev 1.9, American National Standards Institute, Inc.        3. Information regarding FCP I/O devices can be found in Fibre Channel Protocol for SCSI, second version, (FCP-2), American National Standards Institute, Inc.        4. Information regarding the SCSI mapping onto IB can be found in SCSI Remote Direct Memory Access (RDMA) Protocol (SRP), rev 16a, American National Standards Institute, Inc. and        5. Information regarding SRP target port identifiers can be found in Fibre Channel HBA API (FC-HBA), rev 8, American National Standards Institute, Inc.        \nIn an Infiniband configuration containing an I/O Unit (IOU) that enables the host computing system connected to an Infiniband network to access I/O devices, the FCP I/O devices appear to the host as SRP I/O devices. Since all SRP I/O devices are uniquely identified by SRP target port identifiers (IDs), there is a one-to-one correspondence between the worldwide-unique port name (WWPN) for each FCP I/O device and a corresponding SRP target port identifier. SRP target port identifiers consist of two 64-bit quantities. The first is an “IOCGUID” that identifies the IOC that controls the SRP target port. The second is a 64-bit extension. For the case of an IB-to-FC adapter, the extension is set to the WWPN of the FCP I/O device corresponding to the SRP target port identifier. For example, the SRP target port identifier IOCGUID1.WWPNA corresponds to an FCP I/O device with WWPN equal to WWPNA that is accessible through an IOC with an IOCGUID equal to IOCGUID1; SRP target port identifier IOCGUID1.WWPNB corresponds to an FCP I/O device with WWPN equal to WWPNB this is also accessible through an IOC with IOCGUID equal to IOCGUID1; SRP target port identifier IOCGUID2.WWPNC corresponds to an FCP I/O device with WWPN equal to WWPNC that is accessible through an IOC with an IOCGUID equal to IOCGUID2; and so on.\nIn addition to a target port identifier, SRP I/O devices are identified by an IB “service name.” IB service names uniquely identify “service providers,” such as information databases, communication functions, or (as in this case) SRP I/O device functions. The SRP I/O device service name is a 40-byte UTF-8 character string. The first 24-bytes of the service name contains the following character string:                ‘SRP.T10:xxxxxxxxxxxxxxxx’        \nThe “SRP.T10’ part of the service name identifies the service as an SRP target port; the ‘xxxxxxxxxxxxxxxx’ part of the service name is a 16-byte hexadecimal encoding of the 64-bit “extension” portion of the SRP target port identifier. Note that in the case of IB-to-FC adapters, the 64-bit extension of the SRP target port identifier is set to the WWPN of the FCP I/O device corresponding to the SRP target port; therefore, the FCP I/O device corresponding to a particular service name can be deduced by examination of the ‘xxxxxxxxxxxxxxxx’ portion of the service name. The additional 16 bytes of the SRP service name contain null characters. An example of an SRP service name corresponding to an FCP I/O device with WWPN x‘5347 9899 5348 8888’ is ‘SRP.T10:5347989953488888’, followed by null characters.\nSince the service name and SRP target port identifier formats described above contain the FC WWPN of the corresponding FCP I/O device, a host is able to generate the SRP service name for the SRP I/O device that corresponds to the FCP I/O device. Since hosts typically identify FCP I/O devices by their WWPNs, this enables an easy determination of the SRP service name for the corresponding SRP I/O device whenever the host needs to access a particular FCP I/O device with a given WWPN.\nIn order to access the SRP I/O device, however, the host also needs the IB address of the IB-to-FC adapter supporting the service. Prior to this invention, there was no efficient method by which a host could determine the IB address of the IB-to-FC adapter that supports the service name corresponding to the FCP I/O device. Therefore, without the use of this invention, a host attempting to access an FCP I/O device with a particular WWPN would need to poll all of the I/O units (IOUs) in the IB subnet in order to determine those IOUs which are IB-to-FC adapters. The host would then poll each IB-to-FC adapter to determine the particular IB-to-FC adapter (or adapters) that provide access to the FCP I/O device with the desired WWPN. Such a polling operation, which involves sending multiple queries to each IOU in the subnet, is impractical when the IB subnet is of any significant size, and results in unacceptable performance degradation."}
{"text": "Many botanical substances contain chemicals that have been found to be useful for the therapeutic treatment of various medical conditions. Since these chemicals are often present in very small amounts, techniques have been developed to extract these substances and to concentrate the therapeutically active agents. Various methods are available for extraction and purification of such substances, including the use of organic solvents, microwave systems, and supercritical CO2 extraction. Organic solvent-based extractions utilize added solvents that are evaporated to form a concentrated extract, which results in a damp, pasty mass that is typically further spray-dried onto a carrier for delivery. Alternatively, supercritical CO2 extraction is another method of collecting such extracts. This extraction method yields a thick, high viscosity resin, oil, or other fluid-like material that can have a honey-like consistency.\nOne pharmaceutically useful botanical substance is the extract of hops (Humulus lupulus L.). Hops cone flowers contain a variety of active agents, including alpha acids, iso-alpha acids, and beta acids, as well as a number of flavonoids and essential oils. Humulone, one of the alpha acids found in hops, has been demonstrated to suppress cyclooxygenase-2 activity, inhibit angiogenesis, and decrease bone loss. Some other biologically relevant properties of hops constituents include anti-inflammatory, antibacterial, antiviral, antifungal, estrogenic, anti-oxidant, anti-allergenic, anti-carcinogenic, and anti-proliferative properties.\nAs with other botanical substances, dried hops flowers contain very small amounts of alpha acids. Supercritical CO2 extraction and other solvent-based extractions of dried hops cones produce a thick, high-viscosity resin that can contain a high percentage of active hops constituents. While extraction is an effective means of providing alpha and beta acids in a highly concentrated form, the resulting extracts have very low solubility in water. This property can make digesting such extracts difficult, resulting in delayed absorption of the acids and delayed onset of certain therapeutic effects. It would therefore be useful to provide the primary constituents of hops extracts in formulations that are soluble in water. In addition, methods of making such formulations from hops extract resin would be desirable."}
{"text": "The present invention relates to a new and distinct cultivar of rosemary that is grown as an ornamental for containers or landscape. The new cultivar is known botanically as Rosmarinus officinalis and will be referred to hereinafter by the cultivar name ‘ROMAN BEAUTY’.\n‘ROMAN BEAUTY’ was selected by the inventor in Beaulieu, Hants, United Kingdom in 2003 for its improved habit for ease of shipping. The new cultivar ‘ROMAN BEAUTY’ is a result of open pollination among a number of unnamed plants of Rosmarinus officinalis Prostratus Group. The seed parent is a plant of Rosmarinus officinalis Prostratus Group. The pollen parent is unknown.\n‘ROMAN BEAUTY’ is a dwarf perennial ground cover characterized by a compact habit, mid-green foliage, and violet-blue flowers. As an immature, first year plant, the stems of ‘ROMAN BEAUTY’ are faintly tinged with red. As the plant ages and becomes mature, the stems become woody and gray-brown in color.\n‘ROMAN BEAUTY’ differs from its shorter female parent Rosmarinus officinalis Prostratus Group in that ‘ROMAN BEAUTY’ is more upright in habit. Two close comparison plants known to the inventor are Rosmarinus ‘Rampant Boule’ (unpatented) and Rosmarinus ‘Pointe du Raz’ (unpatented). ‘ROMAN BEAUTY’ differs from ‘Rampant Boule’ in that ‘ROMAN BEAUTY’ has a tidier habit; and from ‘Pointe du Raz’ in that ‘ROMAN BEAUTY’ is smaller in habit, and with darker flowers.\nThe first asexual reproduction of ‘ROMAN BEAUTY’ was conducted by the inventor in a cultivated area of Beaulieu, Hants, United Kingdom in 2003. The method used for asexual propagation was vegetative cuttings. Since that time the distinguishing characteristics of ‘ROMAN BEAUTY’ have been determined stable and have reproduced true to type in successive generations of asexual reproduction."}
{"text": "Fuel cells which generate electric current by controllably combining elemental hydrogen and oxygen are well known. One form of a fuel cell consists of an anodic layer, a cathodic layer, and a dense ion conducting electrolyte formed of a ceramic solid oxide. Such a fuel cell is known in the art as a “solid-oxide fuel cell” (SOFC). Hydrogen, either pure or reformed from hydrocarbons, is flowed along the outer surface of the anode and diffuses into the anode. Oxygen, typically from air, is flowed along the outer surface of the cathode and diffuses into the cathode. Each O2 molecule is split and reduced to two O−2 ions catalytically by the cathode. The oxygen ions are conducted through the electrolyte and combine at the anode/electrolyte interface with hydrogen ions to form molecules of water. The anode and the cathode are connected externally through the load to complete the circuit whereby electrons are transferred from the anode to the cathode. When hydrogen for the fuel cell is derived by “reforming” hydrocarbons such as gasoline in the presence of limited oxygen, the “reformate” gas includes CO which is converted to CO2 at the anode. Reformed gasoline and diesel oil are commonly used fuels in automotive fuel cell applications. However, other hydrogen-containing fuels for the reforming process such as, for example, JP8, natural gas, propane, synfuels, alkane alcohols, and coal based fuels may be used as well.\nA single cell is capable of generating a relatively small voltage and wattage, typically between about 0.5 volt and about 1.0 volt, depending upon load, and less than about 2 watts per cm2 of cell surface. Therefore, in practice it is known to stack together, in electrical series, a plurality of cells. The outermost interconnects of the stack define electric terminals, or “current collectors,” which may be connected across a load. A typical prior art SOFC for use as an auxiliary power unit (APU) in a vehicle may comprise about 60 individual fuel cells and may generate, at full power, on the order of 5 kilowatts of electric power.\nA complete SOFC system typically includes auxiliary subsystems for, among other requirements, generating fuel by reforming hydrocarbons as discussed above; tempering the reformate fuel and air entering the stack; providing air to the hydrocarbon reformer; providing air to the cathodes for reaction with hydrogen in the fuel cell stack; and providing air for cooling the fuel cell stack.\nA known shortcoming of a complete SOFC system is that it inherently has a relatively large thermal mass, and consequently, such a system is relatively slow in ramping up to full electric output. Electric output can't begin until the fuel cells are warmed to about 550° C., and a temperature of about 750° C. is required for full output and steady-state operation. The mass of the fuel cell stack along with the induced thermal stresses caused by heating the stack through heated air applied to the cathode dictate the length of time required to produce electricity. Even with known methods for preheating and forced heating of elements in an SOFC system, a prior art SOFC system requires on the order of sixty minutes, starting from ambient temperature, to begin producing usable amounts of electricity. Thus, costly and bulky additional energy storage systems would be required on an APU-powered vehicle to provide power during the fuel cell warm-up period, which causes dissatisfaction to users of such systems. In addition, long start-up times result in a reduction in operating efficiency, especially for intermittent, short-duration operation.\nDuring these start-up periods, unwanted reactions of reformate fuel in the stack is possible, causing deposition of soot (coke) on the relatively cool anode surfaces. Such deposition is undesirable and can result in degradation and eventual failure of the fuel cell stack.\nSteam is a known preventor of carbon (coke soot) forming reactions and a known cleaner of soot from anodes. Therefore, the anodes may be cleaned and prevented from coking by injection of hot steam into the fuel cell stack during warm-up to an operating temperature at which coking does not occur. One approach is to produce steam via vaporizing water from a separate water storage tank, and to then inject the steam into the fuel cell stack along with the reformate. Such a process is undesirable since it requires added complexity of apparatus and logic for storing, supplying, vaporizing, and replenishing water adjacent to the fuel cell system.\nWhat is needed is a means for reducing the start-up period required to bring a large solid-oxide fuel cell system to operating temperature.\nWhat is further needed is a means for preventing coke from depositing on the anodes of a solid-oxide fuel cell during the start-up period required to bring the fuel cell system to operating temperature.\nWhat is still further needed is a means for providing hot steam to an SOFC fell cell stack to remove coke which may form on the anodes while the anodes are cool.\nIt is a principal object of the present invention to reduce the start-up time required for a solid oxide fuel cell system to begin producing electricity.\nIt is a further object of the invention to reduce the need for electrical storage systems in applications wherein a solid oxide fuel cell is an auxiliary electric power unit.\nIt is a still further object of the present invention to minimize coking of anodes during the start-up period required to bring a fuel cell system to operating temperature.\nIt is a still further object of the invention to remove coke deposits that have accumulated on the SOFC anodes."}
{"text": "Micro-feature devices have a large number of very small features that are typically formed in and/or on wafers or other types of workpieces by selectively removing material from the wafer and/or depositing material onto the wafer. For example, features are often formed by (a) constructing a pattern in a layer of resist to form a mask on the wafer, (b) etching holes and/or trenches in the wafer through openings in the mask, and (c) filling the resulting features with dielectric, semiconductive, and/or conductive materials. Photolithographic processes are generally used to transfer the intricate patterns of the features onto discrete areas of the layer of resist.\nA typical photolithographic process includes depositing a layer of radiation-sensitive photoresist material on the wafer, positioning a reticle having a mask pattern over a selected area of the photoresist, and then passing an imaging radiation through the reticle to expose the photoresist in the configuration of the mask pattern. A developer, such as an aqueous base or a solvent, is used to remove either the irradiated areas or the masked areas of the photoresist. For example, when the radiation changes the photoresist from being generally soluble in the developer to generally insoluble, then the developer removes the masked portions of the resist layer. Alternatively, when the radiation changes a photoresist from being generally insoluble in the developer to be generally soluble, then the developer removes the exposed portions of the photoresist.\nExisting lithography processes are capable of creating very complex patterns of extremely small features across the surface of a wafer to form the trenches, vias, holes, implant regions, conductive lines, gates, and other features on a wafer. In a typical application, a lithographic tool transfers the pattern in the reticle to the workpiece by scanning or stepping the pattern across precise areas of the workpiece. As microelectronic devices become more complex, there is a drive to continually decrease the size of the individual features and increase the density of the features across the wafer. This significantly increases the complexity of lithographic processing because it is increasingly difficult to accurately focus the pattern onto the face of the wafer. In many applications, the depth of field for focusing the pattern on the wafer is so small that slight variations in the wafer surface and/or the reticle can adversely affect the quality of the pattern transferred to the wafer.\nOne conventional process to compensate for non-uniformities in reticles is to measure the flatness of the reticles in the photolithography tool before processing the wafers. The topography of the reticles is conventionally measured by detecting light that passes through alignment marks in a perimeter region outside of the pattern area of the reticle. Based on the topographical data of the alignment marks in the perimeter region around the pattern area of the reticle, the topography of the pattern area is estimated. Conventional lithographic tools are then adjusted by tilting the wafer stage and/or adjusting the optics to compensate for variances in the estimated pattern area of the reticle (e.g., the estimated curvature of the reticle).\nOne problem with such conventional processes is that the topography in the pattern area of the reticle is estimated based on the alignment marks in the perimeter region of the reticle. As the feature sizes decrease, this may not provide sufficiently accurate data to compensate for non-uniformities in the pattern area. Moreover, lithographic tools are extremely expensive and it is very costly to use lithographic tools for measuring the non-uniformities in the reticle. Such use of lithographic tools is expected to reduce the throughput of processing wafers because the time period for qualifying the reticles is effectively downtime for processing wafers. Therefore, there exists a need to improve conventional photolithographic processes."}
{"text": "1. Field Of The Invention\nThis invention relates to a thermosetting coating composition and in particular to a thermosetting coating composition of a hydroxyl containing polymer and a polyisocyanate.\n2. Description Of the Prior Art\nThermosetting coating compositions are well known in the art, in particular, thermosetting compositions utilizing polyesters are known, as shown in Sekmakas U.S. Pat. No. 3,457,324, issued July 22, 1969. Also, thermosetting coating compositions of a hydroxyl containing polyester and a blocked polyisocyanate are known as shown in Steinmetz U.S. Pat. No. 3,808,160, issued Apr. 30, 1974. Coating compositions disclosed in these patents form good quality finishes. However, industry requires coatings of improved quality for finishes for appliances such as stoves, refrigerators, washers, dryers and the like, that are highly durable, flexible, scratch resistant, mar resistant, stain resistant and detergent resistant. Also, the industry is in need of a coating composition for coils or sheets of metal that can be applied by conventional coil coating techniques which will provide finishes with improved flexibility and that will withstand post-forming operations that are used to make the aforementioned appliances."}
{"text": "This invention relates to a sanitary napkin for absorption and containment of menstruation discharge.\nThere have been made various attempts to improve fitting of a sanitary napkin to the wearer\"\"s skin by shaping the liquid-absorbent core to be convex in its transversely middle region.\nHowever, adoption of the convex core is accompanied with considerable labor and time, for example, to shape such convex core and then to place the liquid-pervious topsheet closely upon such convex core. In addition, the napkin adopting such convex core becomes bulky in its convex middle region and it is difficult to stack a plurality of napkins into a compact package. Furthermore, the convex region is resistant against a force intending to fold this region and makes it difficult to fold up the napkin as a whole into a conveniently small size.\nAn object of this invention is to provide a sanitary napkin adapted to ensure the same advantageous properties inclusive of a good fitting as the known napkin adopting the convex core can achieve and, in addition, improved to eliminate the problems peculiar to such napkin adopting the convex core.\nAccording to this invention, there is provided a sanitary napkin having a longitudinal direction and a transverse direction intersecting said longitudinal direction, the napkin comprising a liquid-pervious topsheet, a liquid-impervious backsheet and a liquid-absorbent core disposed between the topsheet and the backsheet, a middle region of said napkin in the transverse direction has a rigidity lower than a rigidity of both lateral regions lying adjacent the middle region.\nAccording to one embodiment of this invention, the middle region has a width of 10xcx9c30 mm and extends between longitudinally opposite ends of the absorbent core.\nAccording to another embodiment of this invention, the absorbent core is formed on a lower surface of the middle region with a groove of inverted U- or V-shaped cross-section.\nAccording to still another embodiment of this invention, the absorbent core is bonded to an inner surface of the backsheet over portions of the backsheet underlying the both lateral regions, leaving a portion of the backsheet underlying the middle region free.\nAccording to further another embodiment of this invention, a lower surface of the absorbent core is bonded to the inner surface of the backsheet by means of a plurality of adhesive applied spots extending transversely of the napkin and arranged intermittently in a longitudinal direction thereof.\nAccording to an additional embodiment of this invention, an outer surface of the backsheet except a portion thereof underlying the middle region of the absorbent core is adapted to be adhesively fastened to an undergarment."}
{"text": "The invention relates to a telescopic mechanism, especially for steering columns of motor vehicles, with an internal element, which has a flattening on at least one side, and an external element, which is complementary to the internal element and in which the internal element is guided with roll barrels, which roll at the flattening.\nFor motor vehicles with a positively adjustable steering wheel and/or a safety steering column, the steering column must have a telescopic mechanism, which enables the length of the steering column to be adjusted, so that the position of the steering wheel can be adapted to the respective seated position of the driver and/or so that the steering column can shorten itself in the event of a collision. The internal element and the external element are guided telescopically within one another and together form the steering column. For example, the external element is connected at one end with the steering mechanism, while the internal element carries the steering wheel at the opposite end. The telescopic mechanism is intended to make a smooth axial adjustment of the internal element relative to the external element possible. On the other hand, however, it is to ensure a clearance-free and symmetrical transfer of the steering torque from the internal element to the external element. In this connection, symmetrical means that the twisting of the internal element relative to the external element, which is caused by the torque exerted on the steering wheel, is independent of the direction of action of this torque, so that the same steering sensation is impacted to the driver, no matter which direction the steering wheel is turned."}
{"text": "Sleep is essential for health and quality of life. Insomnia is a growing health problem in the United States. It is believed that more than 10-15 million people suffer from chronic insomnia and up to an additional 70 million people suffer from some form of insomnia each year. Insomnia is a condition characterized by difficulty falling asleep (sleep onset), waking frequently during the night (fragmented sleep), waking too early (premature final awakening), and/or waking up feeling un-refreshed. In the National Sleep Foundation's (NSF) Sleep in America Poll 2005, 42% of survey respondents reported that they awoke frequently during the night, 22% of adults reported waking too early and not being able to return to sleep and 38% reported waking and feeling un-refreshed.\nSleep maintenance difficulty is the most commonly reported symptom in primary care patients with chronic insomnia, and is the most common complaint in depressed patients, medically ill populations, especially those with pain symptoms, and in the elderly.\nMedications commonly used to treat sleep disorders, such as insomnia, include sedative antidepressants, antihistamines, benzodiazepines, and non-benzodiazepine hypnotics.\nAlthough there have been several advances in pharmaceutical treatments for insomnia, it is often hard to find an ideal drug for treating particular forms of insomnia. One common problem is early termination of sleep or premature final awakening. For example, many individuals may wake prematurely and not fall back asleep, thereby failing to achieve a full night of sleep. Many drugs that are effective in inducing or expediting sleep initiation do not provide much effect in maintaining sleep, particularly through the eighth and final hour of sleep period. Drugs that are sufficiently powerful to induce a full eight hours sleep often cause serious hangover effects, i.e., the patient has difficulty awakening and/or feels sedated, sleepy, or disoriented and may demonstrate impairment of psychomotor function.\nIn addition to patients having difficulty with early termination of sleep during the last 60, 90, or 120 minutes of an 8 hour sleep period, other patients have problems with fragmented or disrupted sleep. In other words, those patients awaken one or more times during that time period, then fall asleep again. Such fragmented sleep patterns detract from a feeling of restfulness, and make it less likely that the patient will enjoy restful sleep.\nBoth groups of patients would benefit greatly from a drug that addresses their particular sleep deficiency.\nDoxepin is a tricyclic antidepressant that is known to have beneficial effects in treating insomnia. See, e.g., U.S. Pat. Nos. 5,502,047 and 6,211,229. However, prior to the present invention, doxepin was not known to have particular efficacy in treating premature termination of sleep at the end of an 8 hour sleep period, nor was it known to be efficacious in treating those patients with disturbed sleep patterns during the final 60, 90, or 120 minutes of an 8-hour sleep period. The mean half-life of doxepin is 17 hours, and the half-life of its major active metabolite, desmethyldoxepin, is 51 hours. Thus, when taken at the start of a sleep cycle, a majority of the drug or active metabolite should still be present in the body at the end of the sleep cycle. As a result, it would be expected that dosages of doxepin that are sufficient to address premature final awakenings or last-hour sleep efficiency in the elderly would also cause post-sleep sedation or other undesirable side effects.\nThe present invention describes the surprising ability of doxepin to treat last-hour sleep efficiency and premature final awakenings in patients, without untoward side effects."}
{"text": "This invention relates generally to vehicle air bags and, more particularly, to vehicle air bags having internal tethers for controlling air bag deployment.\nDriver-side air bags have become quite common recently. These air bags are typically stored within the hub of the vehicle steering wheel and deploy to the shape of a flattened sphere. Since the air bag deploys from such a small storage space, it is desirable to assure that the air bag deploys to its full diametral dimension, or as broadly as possible, to cushion the vehicle driver.\nMany driver-side air bags are provided with internal so-called \"suspension belts\", or tethers, to assure full breadth of deployment. These tethers comprise strips of webbing sewn to the air bag about its gas inlet and to the crown of the air bag body opposite the inlet. These tethers are provided to restrain rearward movement of the center of the air bag, thus assuring that the sides of the air bag will deploy.\nSuch arrangements are disclosed in U.S. Pat. No. 4,887,842--Sato and U.S. Pat. No. 5,033,771--Miyauchi et al. The Miyauchi patent is directed to providing a tether attachment which assures that the sewn attachments cannot break so that the center of the air bag cannot deploy further rearward toward the occupant.\nAll air bags for vehicle occupants, whether provided for drivers or passengers, are designed for use in conjunction with conventional seat belt systems. As such, they are normally designed to deploy to engage and cushion a normally seated and belted occupant. However, it has been found that some vehicle occupants fail to utilize the seat belt systems provided. In this instance, it is possible for the deploying air bag to prematurely contact a vehicle occupant who has moved out of the normal seated and belted position.\nIn this instance, it is desirable that the air bag quickly be deployed as broadly as possible to provide contact by the occupant with a full-width air bag. It is also desirable to assure initial downward deployment of the air bag to assure contact with the lower torso of the occupant."}
{"text": "Obesity can be viewed as an energy balance disorder, arising when energy input exceeds energy output, with most of the excess calories converted into triglycerides and stored in the adipose tissue. Medications currently approved for the treatment of obesity attempts to restore energy balance primarily by decreasing energy input by either suppressing appetite or interfering with lipid absorption in the small intestine. Because of the rapid increase in the prevalence of obesity worldwide and the lack of efficacy of current medical therapies, novel pharmacologic therapies for obesity are required.\nOne potential therapeutic strategy involves inhibiting triglyceride synthesis. Although triglycerides are essential for normal physiology, excess triglyceride accumulation results in obesity and, particularly when it occurs in nonadipose tissues, is associated with insulin resistance. DGAT is an enzyme that catalyzes the last step in triacylglycerol biosynthesis. DGAT catalyzes the coupling of a 1,2-diacylglycerol with a fatty acyl-CoA resulting in Coenzyme A and triacylglycerol. Two enzymes that display DGAT activity have been identified: DGAT1 (acyl coA-diacylglycerol acyl transferase 1, see Cases et al, Proc. Natl. Acad. Sci. 95:13018-13023, 1998) and DGAT2 (acyl coA-diacylglycerol acyl transferase 2, see Cases et al, J. Biol. Chem. 276:38870-38876, 2001). DGAT1 and DGAT2 do not share significant protein sequence homology. Importantly, DGAT1 knockout mice are protected from high fat diet-induced weight gain and insulin resistance (Smith et al, Nature Genetics 25:87-90, 2000). The phenotype of the DGAT1 knockout mice suggest that a DGAT1 inhibitor has utility for the treatment of obesity and obesity-associated complications.\nWO2007/126957 discloses a genus of compounds which are disclosed to be inhibitors of DGAT1, and therefore useful in the treatment of a condition or a disorder such as obesity, diabetes and related metabolic disorders. Example 5-23 of said document discloses the compound 2-((1R,4R)-4-(4-(5-(benzo[d]oxazol-2-ylamino)pyridin-2-yl)phenyl)cyclohexyl)acetic acid having the structural formula (I):\n\nNo salt forms of this compound where prepared and the free acid exhibited low intrinsic dissolution and low solubility making large scale manufacture and purification difficult.\nIt is thus important to provide the compound of formula (I) in a salt form and physical form which can be reliably prepared and purified on a large scale, and ideally is stable and does not degrade on storage. The salt form/physical form chosen must also be stable whilst the drug substance is being manufactured as a formulation which is suitable for the intended route of administration chosen. In that respect, it may be necessary to consider physical properties of the salt form in a particular physical form which lead to improved powder handling properties or higher bulk density. In particular, non-hygroscopicity is particularly important in order to obtain good flow characteristics.\nThe properties of the final product should also be predictable and reliably reproducible. For example, material which is obtained in an inconsistent manner, for example, where the water content differs from batch to batch, must be carefully monitored. This leads to added complications in the handling, manufacture, analysis and formulation of the drug substance.\nWhilst one solid state form may exhibit properties which are considered suitable, another form may also have properties which, with the right measures in place, can lead to its successful development into a drug. The decision as to whether a compound is suitable for commercialization thus depends on finding a solid state form of the compound which has the right balance of desirable characteristics."}
{"text": "1. Field of the Invention\nThe present invention relates to a display apparatus and a computer readable medium.\n2. Description of Related Art\nConventionally, a display apparatus for displaying plotted points or approximate graphs has been capable of splitting a display area into two areas, so that plotted points/approximate graphs in different coordinate systems such as an XY coordinate system and an XZ coordinate system can be displayed in each area respectively (e.g., refer to the Japanese Patent Application Laid-Open Publication No. 2003-281102). According to such a display apparatus, it becomes possible to analyze distribution of the plotted points of a three dimensional space simultaneously in a plurality of planes of coordinates and therefore to increase learning effect of a user.\nHowever, by the display apparatus of the above-mentioned patent document, distribution of the plotted points in the three dimensional space and events in the real life cannot be correlated and therefore learning effect of the user is low."}
{"text": "Erythromycins A through D, represented by formula (I) below,\n______________________________________ ##STR3## (I) Erythromycin R.sup.a R.sup.b ______________________________________ A OH CH.sub.3 B H CH.sub. 3 C OH H D H H ______________________________________\nare well-known and potent antibacterial agents. Erythromycin A in particular is widely used to treat and prevent bacterial infection. As with other antibacterials, however, bacterial strains having resistance or insufficient susceptibility to erythromycin have been identified. Another drawback of erythromycins is their poor acid stability, which sometimes results in poor or unpredictable oral absorption. Consequently, numerous investigators have sought to prepare chemical derivatives of erythromycin in an attempt to obtain analogs having modified or improved profiles of antibiotic activity."}
{"text": "Hot melt adhesives are defined as adhesives that are applied from the melt and that build both adhesive and cohesive strength upon resolidification and crystallization of one or more components of the adhesive composition. Because hot melt adhesives are applied without solvents, increasing regulation of VOC emissions in recent years has caused an increase in demand for new and better performing hot melt adhesive formulations.\nHot melt adhesive compositions are formulations that include one or more structural or base polymers and, typically, one or more adjuvants. Adjuvants typically include functional materials or diluents such as tackifying resins, plasticizers, fillers, oils and waxes or other low molecular weight polymers. The base polymer is typically the major component of a hot melt adhesive formulation and contributes cohesive strength to the adhesive. Copolymers of ethylene and a mono-olefinically unsaturated polar comonomer such as vinyl acetate, methyl acrylate, acrylic acid and the like, are widely used as the base polymer component in many hot melt adhesive compositions. Another class of hot melt adhesives employs copolymers of ethylene or propylene with one or more relatively nonpolar α-olefins such as 1-butene, 1-hexene, and the like. Hot melt adhesive formulations employing nonpolar polymers, while less common in the industry than their EVA and other similar polar polymer based counterparts, are known. Examples of some such formulations are disclosed in U.S. Pat. Nos. 4,072,735; 5,118,762; and 4,568,713 to name just a few examples.\nU.S. Pat. No. 6,747,114 discloses hot melt adhesive formulations employing metallocene polymerized propylene copolymers, wherein the comonomer is ethylene and/or a C4-C20 α-olefin as the base polymer in the formulation. The highly isotactic or syndiotactic propylene content of the metallocene polymerized base polymers imparts crystallizable content to the base polymers; the metallocene polymerized propylene copolymers are called semicrystalline copolymers. As used herein and as applied to semicrystalline copolymers, the term “crystallizable” describes those sequences which are mainly amorphous in the undeformed state, but can crystallize upon stretching, annealing or in the presence of a nucleating agent, such as a crystalline compound or another segment within the polymer. Crystalline content of the solidified semicrystalline copolymers increases the cohesive strength of the hot melt adhesives. Hot melt adhesive formulations based on metallocene polymerized propylene-based semicrystalline copolymers can eventually build sufficient crystalline content over time to achieve good cohesive strength in the formulation.\nHowever, such formulations suffer from the drawback that the slow rate of crystalline content formation, retards cohesive strength formation. For many potential end uses, the rate of crystallization of the crystallizable content of the metallocene polymerized propylene-based semicrystalline copolymers is unacceptably long and results in a significant delay in effective adhesive bonding and reduced productivity. Attempts to accelerate the rate of crystallization have centered on the use of nucleating agents. U.S. Pat. No. 6,747,114 discloses metallocene polymerized propylene-based semicrystalline copolymers hot melt adhesive formulations having about 10 to 25 wt. % of a Fischer-Tropsch wax and/or 5% of a low molecular weight isotactic polypropylene homopolymer in order to affect crystalline content in the semicrystalline copolymer-based hot melt adhesive.\nWhile the nucleation agents described above effectively increase the ultimate degree of crystalline content, as evidenced e.g. by a higher ultimate tensile strength of the adhesive, they do not address the rate of nucleation and crystal formation at a given temperature. Slow rate of crystallization is evidenced by a long effective set time: that is, a relatively long period of time between application of the hot melt adhesive formulation and the ability of the formulation to cause a bond to form between two adherends of sufficient strength for the envisioned application. In fact, the effective set time of the metallocene polymerized propylene-based semicrystalline copolymer hot melt adhesive formulations of the art is too slow for many industrially significant commercial applications.\nAdditionally, the high level of wax included in some such formulations can result in lower ultimate adhesive strength when the adhesive is completely crystallized. For example, one known application of Fischer-Tropsch waxes is as a mold release. The greater the amount of wax in the formulation, the greater the amount present at the adhesive-substrate interface and the greater its effect upon reduced ultimate adhesion and release.\nThus there is a need in the industry to obtain hot melt adhesive compositions based on metallocene polymerized propylene-based semicrystalline copolymer that builds crystalline content quickly. There is a need in the industry to obtain such adhesives having a low wax content. There is a need for such hot melt adhesive compositions to have the properties of melt rheology, adhesion, and cohesive strength widely associated with conventional hot melt adhesive formulations."}
{"text": "Portable electronic devices continue to get smaller and incorporate more functions, such as traditional personal digital assistant (“PDA”) functionality with cellular telephony and wireless email capabilities. In addition to functions oriented toward the business user, it is also known to incorporate music and video players as well as camera applications for consumer market devices.\nConventional film cameras use a photosensitive film to capture an image, whereas digital cameras use electronic photosensors such as charge coupled devices (CCDs) or complimentary metal oxide semiconductor (CMOS) chips. The term “photosensor” as used in this specification means any device(s) or material(s) capable of receiving and capturing radiant energy, and being at least partially capable of converting the radiant energy into electronic signals that become a virtual representation of the optical image. A CCD or CMOS “camera-on-a-chip” includes an array of very fine electronic “picture elements” or “pixels” arranged in horizontal rows and vertical columns that define an image resolution matrix.\nU.S. Pat. No. 5,841,126 describes an exemplary camera chip that may be incorporated into a portable electronic device.\nOne problem associated with such photosensor arrays is the introduction of an image artifact referred to as “flicker” when the camera is capturing an imaged scene that is illuminated by a fluorescent light source. Flicker occurs as a result of periodic variations of light intensity in correspondence with the frequency of the alternating current (AC) power that interfere with the arrays' ability to capture all of the image information used to form an image frame.\nA fluorescent lighting system that is powered by a source of 60 Hz alternating current will exhibit periodic peaks of intensity at a rate of 120 Hz, i.e., twice the frequency of the alternating current. However, in many European countries, the AC power waveform has a frequency of 50 Hz, so that the flicker frequency of concern is 100 Hz. Thus, unless the photosensor chip includes a mechanism for addressing “beats” at this frequency, an image of a gray background captured under illumination by fluorescent lighting will include readily apparent amplitude modulations of the light intensity in a particular direction (typically the vertical direction), since the light level will vary with the capture of different lines of the image.\nA number of solutions have been employed to eliminate these “beats.” These include filtering systems that filter out the beat frequency, phase locking systems that attempt to lock on to the 100 Hz intensity peaks and synchronize frame capture, and a variety of other techniques.\nOne such other technique is set forth in U.S. Pat. No. 6,271,884 to Chung et al., which describes a digital camera with constant frame rate, but with an adjustable integration time. The integration time is defined as the amount of time that a particular sensor is permitted to capture light energy for each frame. When the camera is used in an environment having 60 Hz fluorescent lighting, the integration time is set at a multiple of 8.33 milliseconds, whereas the integration time is changed to a multiple of 10 milliseconds when the environment utilizes 50 Hz fluorescent lighting. A number of options are set forth for between the two integration times. One suggestion is that the system could detect the country in which it is operating based on system configuration data, although absolutely no details are provided on how this would be implemented.\nOther approaches are known in the art for detecting flicker, such as U.S. Pat. No. 7,187,405, which detects specific repeating patterns of signal variations by processing columnar information from the device's two-dimensional sensor array.\nIt is also known in the art to select camera printing options for printing parameters such as printed image size and paper size in accordance with location information derived from a location subsystem, such as a GPS receiver, as described in U.S. Pat. No. 7,126,639.\nIt is contemplated that improvement is possible over the above-described prior art, by simplifying the manner of detecting location of a digital camera for the purpose of setting the preferred fluorescent operating mode to eliminate flicker."}
{"text": "The clinical objective in treatment of peptic ulcer disease is to decrease gastric acid secretion, based on the principle \"no acid, no ulcer.\" Traditional peptic ulcer disease therapy involves control of diet and the use of antacids and anticholinergics.\nThere is evidence indicating that histamine may be the final common pathway for stimulation of gastric secretion. This effect of histamine is mediated via H.sub.2 -receptors and is not inhibited by the classical antihistamines, which are H.sub.1 -receptor blockers. A number of specific H.sub.2 -receptor blocking agents (H.sub.2 -receptor antagonists) are now known. These compounds inhibit basal acid secretion, as well as secretion by other known gastric acid stimulants, and are useful in the treatment of peptic ulcers.\nBurimamide (IIa) was the first clinically effective H.sub.2 -receptor antagonist. It inhibits gastric secretion in animals and man, but oral absorption is poor.\n______________________________________ ##STR2## II IIa; R.sup.4H, ZCH.sub.2, XS Burimamide b; R.sup.4CH.sub.3, ZS, XS Metiamide c; R.sup.4CH.sub.3, ZS, XNCN Cimetidine ______________________________________\nMetiamide (IIb), a subsequently evaluated H.sub.2 antagonist, is more potent than burimamide and is orally active in man. Clinical utility was limited, however, owing to toxicity (agranulocytosis). Cimetidine (IIc) is as effective an H.sub.2 antagonist as metiamide, without producing agranulocytosis, and has recently been marketed as an anti-ulcer drug. The half-life of cimetidine is relatively short, thereby necessitating a therapeutic regimen of multi daily doses of 200-300 mg. tablets. There is thus a need for anti-ulcer agents which are longer acting and/or more potent than cimetidine.\nReviews on the development of H.sub.2 antagonists, including those discussed in the preceding paragraph, may be found in C. R. Ganellin, et al., Federation Proceedings, 35, 1924 (1976), in Drugs of the Future, 1, 13 (1976), and in references cited therein. Relevant patents are as follows:\nBelgian Pat. No. 841,814 (Farmdoc 90568X) discloses inhibitors of histamine-stimulated gastric secretion having the formula ##STR3## in which HET is one of eight named heterocyclic rings (including pyridyl) which may be substituted by (lower)alkyl, hydroxyl, amino or halogen; Z is sulfur or CH.sub.2 ; X is S, CHNO.sub.2, NCN or NH; Y is NH.sub.2, (lower)alkylamino, di(lower)alkylamino, (lower)alkoxy, phenylethyl, imidazolylethyl, allyl, trifluoroethyl or (CH.sub.2).sub.n R; n is 1-12; and R is OH, (lower)alkoxy, NH.sub.2 or (lower)alkylamino; provided that, when X is NH, Y is trifluoroethyl or (CH.sub.2).sub.n R; and when X is NCN, Y may not be amino or (lower)alkylamino.\nBelgian Pat. No. 804,144 (Farmdoc 19437V) discloses inhibitors of histamine-stimulated gastric acid secretion having the formula ##STR4## in which HET is a 5 or 6 membered heterocyclic ring containing nitrogen (pyridine is named), which may be substituted by alkyl, halogen, CF.sub.3, OH or NH.sub.2 ; m and n are each 0-4 and the sum of m and n is from 2 to 4; Z is sulfur, oxygen, NH or CH.sub.2 ; and R.sub.1 is hydrogen or (lower)alkyl.\nU.K. Pat. No. 1,421,792 discloses H.sub.2 -receptor inhibitors of the formula ##STR5## wherein X and Y, which may be the same or different, are hydrogen, nitro, cyano or SO.sub.2 Ar, but may not both be hydrogen; R is hydrogen, (lower)alkyl or Het(CH.sub.2).sub.m Z(CH.sub.2).sub.n ; Z is sulfur or methylene; m is 0, 1 or 2 and n is 2 or 3 provided that the sum of m and n is 3 or 4; Het is an imidazole, pyridine, thiazole, isothiazole, oxazole, isoxazole, triazole or thiadiazole ring which is optionally substituted by (lower)alkyl, hydroxy, halogen or amino; and Ar is phenyl, optionally substituted by halogen, methyl or amino.\nBelgian Pat. No. 844,504 (Farmdoc 07558Y) discloses H.sub.2 -receptor inhibitors of the formula ##STR6## wherein R.sub.1 is (lower)alkyl or --(CH.sub.2).sub.p A; p is 2-4; A is hydroxy, (lower)alkoxy or dimethylamino; R.sub.2 is hydrogen, (lower)alkyl, (lower)alkoxy, amino, methylamino or halogen; or --OR.sub.1 and R.sub.2 taken together form the group --O(CH.sub.2).sub.q O-- joined to two adjacent carbon atoms of the pyridine ring; q is 1-4; n is 2 or 3; Z is sulfur or --CH.sub.2 --; X is sulfur, CHNO.sub.2, NH, NCN or NOH; and Y is inter alia hydrogen, (lower)alkyl or 2-hydroxyethyl.\nU.S. Pat. No. 4,112,234 discloses histamine H.sub.2 -receptor inhibitors of the formula ##STR7## wherein R.sup.1 is a straight or branched chain alkynyl group containing from 3 to 9 carbon atoms, and processes for the preparation thereof.\nU.S. Pat. No. 4,060,621 discloses histamine H.sub.2 -receptor inhibitors of the formula ##STR8## wherein Het is a pyridyl ring optionally substituted by (lower)alkyl, trifluoromethyl, hydroxyl, halogen or amino; Z is sulfur, oxygen, NH or a methylene group; m is 0, 1 or 2 and n is 2 or 3, the sum of m and n being from 2 to 4; X is COR.sub.3, CSR.sub.3, SO.sub.2 R.sub.4, NCHR.sub.5 or, when Z is methylene, may be nitro; R.sub.3 is (lower)alkyl, (lower)alkoxy or, when Z is sulfur, oxygen or NH, may be amino; R.sub.4 is (lower)alkyl, trifluoromethyl, amino or substituted or unsubstituted aryl, such as phenyl optionally substituted by halogen, (lower)alkyl or amino; R.sub.5 is substituted or unsubstituted aryl, such as phenyl; and R.sub.1 is hydrogen or (lower)alkyl such as methyl; and pharmaceutically acceptable salts thereof. U.S. Pat. No. 3,971,786 contains a substantially identical disclosure."}
{"text": "Systems that network voicemail systems to facilitate a distributed office, include, without limitation, networks by Cisco®, Advanced Voice®, etc. Networks like that of Cisco® have included voice over Internet protocol (VoIP) communications. However, prior networked voicemail systems do not transfer the session of the caller from one system to another system of the network using speech VoIP to transfer the session. Instead, prior systems transfer messages or calls but not the entire session. Accordingly, since the entire session is not transferred, the caller typically cannot speak to several people at the distant location, leave messages for several people at the location, broadcast messages at the location, etc. all as though they were at the other location.\nIn view of the above, there is a continuing need for a network of attendant, messaging and content-based speech systems to facilitate a distributed office wherein each system can transfer a session of a caller to another system at a distant location so that the caller can speak to several people at the distant location, leave messages for several people at the location, broadcast messages at the location, etc. all as though they were at the other location.\nAdditionally, there is a continuing need for a network of self-synchronizing attendant, messaging and content-based speech systems that automatically updates other systems in the network with extension database additions, deletions or other changes.\nAs will be seen more fully below, the present invention is substantially different in structure, methodology and approach from that of the prior e-trading systems and methods."}
{"text": "1. Technical Field\nThis invention relates to electric motors, and more particularly to a power tool having an electric motor adapted to be powered by either an AC power source or a DC battery.\n2. Discussion\nElectric motors are used in a wide variety of applications. In the power tool industry, electric motors are used to power various tools such as drills, grinders, circular saws and other various portable, electrically powered tools. In recent years, battery powered electric hand tools have become increasingly popular. The provision of the battery has greatly enhanced the utility of such hand tools, allowing their use in applications where an AC power source is unavailable.\nIn spite of the advantages of having a battery powered electric hand tool, it would nevertheless still be desirable to be able to use AC power in the event the battery of the power tool becomes discharged to a point where the tool can no longer be used, but an AC power source is readily available. In such situations, simply connecting the electric power tool to the AC power source via an electric extension cord would enable the power tool to be used to finish the task at hand.\nAccordingly, it is a principal object of the present invention to provide a power tool capable of being powered by either an AC or a DC power source without requiring modification to the tool.\nIt is a further object of the present invention to provide a power tool having a motor which is capable of being powered by an AC or a DC power source, and where the motor does not require the overall dimensions of the power tool to be increased significantly, does not increase the overall weight of the power tool significantly, and further does not significantly increase the cost of manufacturing the power tool."}
{"text": "1. Field of the Invention\nThe present invention relates to position measurement methods, position control methods, measurement methods, loading methods, exposure methods and exposure apparatus, and device manufacturing methods, and more particularly to a position measurement method in which position information of a plate mounted freely detachable on a moving body is measured, a position control method that uses the position measurement method, a measurement method in which information related to a plate where an opening is formed in order to mount an object is measured, a loading method that uses the measurement method, an exposure method that utilizes the loading method and an exposure apparatus suitable for performing each of the methods described above, and a device manufacturing method that uses the exposure apparatus or the exposure method.\n2. Description of the Related Art\nConventionally, in a lithography process for manufacturing electronic devices such as a semiconductor device (an integrated circuit or the like), a liquid crystal display device, or the like, a reduction projection exposure apparatus (the so-called stepper) by the step-and-repeat method that transfers a pattern formed on a mask or a reticle (hereinafter generally referred to as a ‘reticle’) onto a photosensitive object such as a wafer or a glass plate (hereinafter generally referred to as a ‘wafer’) on which a resist (a photosensitive agent) is coated, or a projection exposure apparatus (the so-called scanning stepper) by the step-and-scan method is mainly used.\nDue to higher integration and finer circuit patterns of the semiconductor devices, in order to improve the resolution of the projection optical system equipped in the projection exposure apparatus, the wavelength of the exposure light (exposure wavelength) is becoming shorter while the numerical aperture (NA) of the projection optical system is gradually increasing. Meanwhile, depth of focus is becoming smaller, due to such shorter exposure wavelength and increasing numerical aperture (larger NA). The exposure wavelength is presumed to be shorter in the future, and if such a situation continues, the depth of focus may become so small that margin shortage may occur during the exposure operation.\nTherefore, as a method of substantially shortening the exposure wavelength while increasing (widening) the depth of focus when compared with the depth of focus in the air, the exposure apparatus that utilizes the immersion exposure method is beginning to gather attention. As the exposure apparatus that utilizes the immersion method, an exposure apparatus that performs exposure in a state where the space between the lower surface of the projection optical system and the wafer surface is locally filled with liquid such as water or an organic solvent is known (refer to, for example, the pamphlet of International Publication No. WO99/49504). In the exposure apparatus according to the pamphlet, the resolution is improved utilizing the fact that the wavelength of the exposure light in the liquid becomes 1/n of the wavelength in the air (n is the refractive index of the liquid which is normally around 1.2 to 1.6), and also the depth of focus is substantially increased n times when compared with the case where the same resolution is obtained by a projection optical system (supposing that such a projection optical system can be made) that does not employ the immersion method. That is, the depth of focus can be substantially increased n times than in the air.\nRecently, in wafer stages of the exposure apparatus, a proposal has been made of disposing a freely detachable plate that forms a flat section substantially flush with the wafer in the periphery of the wafer held by the wafer stage. In the case of using such a detachable plate in the wafer stage, the position of the plate has to be precisely known.\nIn addition, in the case of using such plate in the wafer stage, an opening (such as an opening with a circular shape in the case of a semiconductor wafer) for positioning the wafer has to be formed in the center of the plate. However, for example, in the case the degree of roundness of the circular opening of the plate is low and the circular shape is deformed or in an ellipse, the gap between the circumferential surface of the wafer and the inner circumference surface of the opening becomes uneven, and inconveniences could occur, such as the wafer coming into contact with the inner wall surface of the opening of the plate, or not being able to insert the wafer into the opening of the plate.\nIn addition, because the gap between the inner wall surface of the opening of the plate and the wafer is extremely small, smooth loading operation of the wafer will be difficult if the relative position of the wafer and the plate is not accurately aligned when loading the wafer.\nIn addition, in the case of the exposure apparatus that utilizes the immersion method, there was the risk of the liquid flowing into parts where the gap between the inner circumference edge of the opening of the plate and the outer circumferential edge of the wafer is large."}
{"text": "1. Field of the Invention\nThis invention relates generally to a leader for a magnetic or other recording tape, and more particularly is directed to an improved leader for a magnetic or other recording tape which, when secured to one or both of the ends of the latter, may be detected to indicate the presence of the adjacent end of the recording tape in an audio or video recording and/or reproducing apparatus or the like.\n2. Description of the Prior Art\nIn tape cassettes for use in audio or video recording and/or reproducing apparatus, it is known to provide leaders connected between the opposite ends of a magnetic recording tape and the hubs of the supply and takeup reels on which the tape is wound so that one or the other of the tape leaders can be detected to indicate the full unwinding of the tape from a respective reel during a recording, reproducing, fast-forward or rewind operation of the apparatus in which the tape cassette is employed. One existing type of tape leader is formed of a flexible base of polyester film having a thin layer or film of aluminum or copper provided on one surface of the base, as by vapor depositing or plating the metal on the base or by adhering or bonding a metal film to the base. A tape leader of the foregoing type can be conveniently detected by a resonance sensor including an oscillator having a L-C resonance circuit, a detector for detecting the output of the oscillator, and a warning and/or shut-off mechanism which responds to detection of a change in the frequency of the oscillator output or a halt in such output. At least the inductance element of the L-C resonance circuit is included in a sensor which is disposed adjacent the path of the recording tape so that, when the tape leader with an aluminum or copper film or coating thereon is disposed adjacent the sensor, the inductance or sharpness of the resonance circuit is changed or decreased and, accordingly, the oscillating frequency is changed or the oscillation is halted. For example, in an existing video tape recorder (VTR), when the tape leader is detected, as aforesaid, an automatic shut-off mechanism is actuated or energized to halt the recording, reproducing, fast-forward or rewinding operation and to eject the tape cassette from the VTR. Further, in such VTR, detection of the tape leader causes energizing of a warning lamp and/or buzzer or other audible signal.\nHowever, the existing tape leader, as described above, has a relatively poor ability to return to its original flat condition after being bent. Therefore, after frequent use of the recording tape having such leader or leaders, the latter may become crumpled or snap, or cracks or fissures may develop in the conductive metal film or coating thereof. The foregoing problems are especially encountered when the recording tape with a leader or leaders is made to follow a sinuous path about numerous tape guides in the tape cassette and/or in the recording and/or reproducing apparatus with which the cassette is used. When a tape leader is damaged, as aforesaid, it is difficult to replace the same, particularly if the associated recording tape and leaders are incorporated in a tape cassette. Furthermore, when the tape leader or leaders become crumpled, the crumpled condition of the leader material may injure or cause uneven winding of the magnetic or other recording tape about the crumpled leader on the associated supply or takeup reel with the result that dropout of the video or other recorded signals may be experienced.\nMoreover, in the above described existing tape leader, the aluminum or copper film is corroded or oxidized when the recording tape is exposed to a humid atmosphere for an extended length of time. Such corrosion or oxidation of the aluminum or copper film of the tape leader changes its electric conductivity, with the result that the sensing of the tape-end is undesirably influenced, that is, the sensor can no longer reliably detect the tape leader when the latter is adjacent thereto."}
{"text": "It has long been known in the prior art to employ a plug to effect a blockage of a normally open passage, when for some particular reason such action is deemed necessary. Moreover, the prior art is known to be replete with examples of various forms of plugs, which differ one from another in terms of configuration, mode of construction, etc. On the other hand, even though a multiplicity of plugs of different designs exist in the prior art, there continue to arise applications wherein it would be desirable to be able to employ a plug. However, due to the peculiar nature of the application, the plugs of known design are unsuitable for use therein.\nOne such application involves the accomplishment of the blockage of the inlet nozzle and/or the outlet nozzle of a steam generator, the latter being of the type that is suitable for employment in cooperative association with a nuclear reactor to form therewith a component part of a nuclear power generating system.\nPeriodically, there arises a need to conduct maintenance on and/or to effect minor repairs of the internal components of the steam generator. To accomplish such tasks, it is necessary for one or more persons to physically enter the steam generator. Ingress and egress from the steam generator by such persons is accomplished through suitable means such as for example an access port. While such persons are working in the steam generator, it is desirable that a blockage of the inlet and the outlet nozzle of the steam generator be effected to insure that there will be no fluid flow through the steam generator, which might imperil the lives of the people working therein.\nOne such design of a nozzle plug for a nuclear steam generator is disclosed in co-pending application Ser. No. 349,565 by Wentzell filed Feb. 17, 1982, entitled \"Nozzle Plug For Submersible Vessel\", which is herein incorporated by reference. This prior design is substantially comprised of a plurality of plate sections which are assembled within the steam generator to form the plug which is secured within the particular nozzle to be blocked. As this assembly and securing of the nozzle plug within the steam generator must be accomplished in an area where radioactivity may be present, it is desirable to be able to perform these operations as quickly and as efficiently as possible."}
{"text": "The present invention relates to processes and apparatuses for preventing food-borne illnesses by pasteurizing the surfaces of food products.\nIn spite of increasingly stricter government regulations and ongoing efforts in the meat, fish and poultry industries, outbreaks of illnesses caused by food-borne bacteria continue to occur on a regular basis. In 1999 alone, 25 people in the United States died as a result of consuming contaminated meat or poultry products. The presence of any illness-causing bacteria on ready-to-eat (RTE) and other precooked meat, poultry, and fish products (e.g., sliced ham, beef, or turkey logs, Virginia hams, oven-roasted turkey, spiral honey-baked hams, etc.) is of particular concern because these products typically are not recooked or sufficiently reheated prior to consumption.\nThe threat posed to public health and safety by food-borne pathogens is extremely wide-spread and the economic impact of even a single outbreak can be staggering. Recently, one company issued a voluntary recall of approximately 16.7 million pounds of RTE turkey and chicken products because of possible contamination problems at just one of the company\"\"s processing facilities. It is reported that possible contamination problems were traced to the facility following the occurrence of certain illnesses identified by state health departments. The recalled products had been distributed nationwide and to some foreign countries.\nExamples of particularly serious bacterial contaminants include listeria, salmonella, and E-coli. These contaminants are especially problematic when dealing with RTE and other precooked meat, poultry and fish products and are often present in processing environments in spite of diligent efforts to eliminate them. Listeria monocytogenes, for example, is tolerant to salt, can grow at low temperature, has a high heat tolerance, and can form stubborn biofilms on processing surfaces.\nWhen dealing with RTE and other precooked products, the risk of contamination is particularly high after cooking and prior to completing the final packaging process. Such meat, poultry, and fish products are typically cooked to an internal temperature of at least 160xc2x0 F. so that, at the end of the cooking process, no listeria, salmonella, or E-coli should be present. However, prior to packaging, the cooked product will typically be chilled or otherwise allowed to cool to an internal temperature of about 40xc2x0 F. or less and can be contaminated by airborne pathogens and/or bacteria present on conveyor surfaces, on processing equipment, in condensate drippage, or on the hands of workers in the processing area.\nIn response to continuing outbreaks of food-borne illnesses, governmental regulatory agencies continue to impose stricter regulations. USDA and FDA authorities have now established zero tolerance requirements for listeria monocytogenes and salmonella in all RTE food products.\nUnfortunately, as evidenced by the continuing occurrence of illnesses and deaths from food-borne bacteria, the current industry practices and procedures for dealing with these problems are not sufficiently reliable and are inadequate to meet the zero tolerance requirements now imposed by regulatory agencies. Current procedures commonly consist of maintaining the cooked product in a xe2x80x9ccleanxe2x80x9d environment (a xe2x80x9cclean roomxe2x80x9d) and conveying the product along a critical control path designed to prevent bacteria from reaching the product surfaces. In an effort to keep the processing environment clean, the air within the clean environment is typically filtered and operators are instructed to wash all contact surfaces between shifts and to wash their hands each time they enter the clean area.\nIt is thus apparent that a need presently exists for a process which will consistently and effectively kill surface bacteria present on food products, particularly on RTE and other precooked meat, poultry, and fish products, and will meet and exceed all governmental regulatory requirements. A need particularly exists for such a process which will not alter the surface characteristics or internal characteristics of the products in any significant way.\nThe present invention provides a surface pasteurizing system which satisfies the needs and alleviates the problems discussed above. The inventive system can be used for pasteurizing the surface of generally any raw or cooked food product and is particularly well suited for treating precooked whole muscle, emulsified, or other meat, poultry and fish products. In its most preferred embodiments, the inventive system is effective for destroying bacteria without producing any substantial change in the color or other characteristics of the product.\nIn one aspect, the present invention provides a process for preventing food-borne illness comprising the step of heating a surface of a food product in a manner effective to achieve at least a 3 log reduction in live bacteria on the surface without causing any substantial color change (i.e., any change in color readily discernable by the naked eye) in the surface.\nIn another aspect, the present invention provides a process for preventing food-borne illness comprising the step of heating a surface of a food product to a temperature of at least 160xc2x0 F., the step of heating being conducted in a manner effective such that no substantial increase in the internal core temperature of the food product (i.e., no increase amounting to as much as 1xc2x0 F. or more) occurs.\nIn yet another aspect, the present invention provides a process for preventing food-borne illness comprising the steps of (a) continuously conveying a precooked food product through a continuous infrared oven at an operating temperature of at least 500xc2x0 F. and (b) heating a surface of the precooked food product in the infrared oven for a time sufficient to bring the surface to a temperature of at least 160xc2x0 F. The food product is selected from the group consisting of meat, poultry, and fish products.\nFurther objects, features, and advantages of the present invention will be apparent to those skilled in the art upon examining the accompanying drawings and upon reading the following description of the preferred embodiments."}
{"text": "The eXtensible Markup Language (XML) is a markup language that allows users to describe data in hierarchically-structured documents or equivalent files. In general, the data is not only present in an XML document, but is described in some way. For example, various sets of text in an XML document might be tagged as separate paragraphs, whereby a program interpreting the document would know something about the text's organization.\nXML is a simplified subset of SGML (Standard Generalized Markup Language) that removes some of SGML's more complex features to simplify programming. XML is a defined non-proprietary standard, so XML-formatted information is accessible and reusable by any XML-compatible software, in contrast to proprietary formats used by many conventional programs such as traditional word processors. In other words, XML can be used to store any kind of structured information in a manner that enables it to be communicated between computers, including those that are otherwise unable to communicate. The format is robust, persistable and verifiable.\nXML allows the flexible development of user-defined document types that are stored, transmitted and/or processed in some manner, while providing information content that is richer and easy to use, (e.g., relative to HTML), because the descriptive and hypertext linking abilities of XML are much greater than those of HTML.\nAs XML and XML documents are becoming extremely popular, various tools are needed to work with XML technology. One such tool that would benefit users would provide a way to compare two XML documents. File comparison has a wide range of uses, generally known from word processor utilities and the like that perform line-oriented comparisons, such as those that compare text.\nHowever, while such line-oriented comparisons systems are straightforward to implement, they are also rather limited, and do not fit the hierarchical nature of the structure of XML documents. What it needed is a comparison method and system that are tree-oriented, to match the hierarchical structure of structured documents such as XML documents."}
{"text": "1. Field of the Invention\nThe present invention relates to an optically active compound and a liquid crystal composition containing said compound The compound and composition containing said compound, of the present invention show a ferroelectric liquid crystal phase, and accordingly are useful as an electrooptic switching element such as liquid crystal display device or the like and can be used in liquid crystal optical modulators.\n2. Discussion on the Related Arts\nLiquid crystal display devices have various excellent features such as low-voltage operability, low power consumption, being thin and light-weight, and easy on the eye to be a light-receiving type. Accordingly, they are in wide use as various display devices.\nLiquid crystal display devices using a nematic liquid crystal called a twisted nematic mode (TN mode) are in use currently. However, display devices using this nematic liquid crystal have had a drawback of being very slow in response as compared to luminescent type display devices such as CRT, EL and the like. Consequently, when such display devices are applied in a display device, particularly a large scale display device capable of displaying a large amount of information, it is impossible to obtain a display of good contrast. Thus, the liquid crystal display devices using a nematic liquid crystal have had a limitation for wide applications. There has recently been developed a liquid crystal device using a nematic liquid crystal called a super twisted nematic type (STN mode) or SBE and capable of giving a display of improved contrast. Even in this STN type liquid crystal display device, however, the response is not sufficient and therefore said device finds a limitation in application to a display capable of displaying a still larger amount of information. Hence, various attempts are under way to develop a new liquid crystal display system giving an excellent response.\nFerroelectric liquid crystals have a memory function and give a high speed response, and accordingly their application to a large scale display is highly expected. As liquid crystals having ferroelectric properties, there are known those showing a chiral smectic C phase, a chiral smectic H phase, a chiral smectic J phase, etc. Of these ferroelectric liquid crystals, those showing a chiral smectic C phase are thought to have highest practical usability.\nFerroelectric liquid crystals showing a chiral smectic C phase were first synthesized in 1975 by R. B. Meyer et al., and their typical examples include 2-methylbutyl 4-(4'-n-decyloxybenzylideneamino)cinnamate (hereinafter abbreviated to DOBAMBC) [J. Physique, 36, L-69 (1975)].\nA thin film liquid crystal cell was prepared response of .mu. sec order [N. A. Clark et al., Appl. Phys. Lett., 36, 89 (1980)]. Since that time, there was started the development of optical modulation devices (e.g. liquid crystal device, photo-printer head) using a ferroelectric liquid crystal showing a chiral smectic C phase (hereinafter may be referred to simply as \"ferroelectric liquid crystal\").\nAs a result, a number of ferroelectric liquid crystal compounds showing a chiral smectic C phase have been developed since then and various ferroelectric liquid crystal compounds are known currently. However, no ferroelectric liquid crystal compound is found yet which has satisfactory reliability and capability for use in a display device, particularly a large scale one, etc.\nIn order for a ferroelectric liquid crystal to be practically used as a liquid crystal display device, etc., the liquid crystal must be superior in high speed response, orientation, memory function, threshold, temperature dependences of these properties, etc. Also, the ferroelectric liquid crystal is required to show a chiral smectic C phase over a wide temperature range so that it can operate at a sufficiently wide temperature range including room temperature and further to have excellent physical and chemical stabilities.\nIn order for a ferroelectric liquid crystal to have, in particular, excellent physical and chemical stabilities, good high speed response and good memory function, the liquid crystal must have a large spontaneous polarization.\nAmong the so far developed ferroelectric liquid crystal, no compound is found yet which satisfies the above requirements. For example, the above mentioned DOBAMBC, being a liquid crystal of Schiff's base type, is insufficient in chemical stability to water, light, etc. and moreover, has a small spontaneous polarization of 4 nC/cm.sup.2 or below.\nEster type liquid crystals are reported as a ferroelectric liquid crystal which is chemically stable. However, these liquid crystals are not satisfactory because they have no sufficiently large spontaneous polarization and no sufficiently wide temperature range of chiral smectic C phase.\nIn order to obtain a large spontaneous polarization, there were synthesized those compounds having two asymmetric carbon atoms as an optically active group essential for the expression of chiral smectic C phase.\nThese compounds include, for example, liquid crystal compounds having a dichiral epoxide side chain [David M. Walba et al., Journal of American Chemical Society, 108, 7424 (1986)] and liquid crystal compounds having a halogen atom and a methyl group on two adjacent asymmetric carbon atoms [Japanese Patent Application Kokai (Laid-Open) Nos. 168780/1985, 218358/1985, 68449/1986, 40/1987, 46/1987, 103043/1987, 111950/1987, 142131/1987, 175443/1987, etc.].\nAs a typical example of the above liquid crystal compounds, there is (s)-3-methyl-2-chlorobutyl 4-(4'-octylcarbonyloxy) biphenylcarboxylate [Japanese Patent Application Kokai (Laid-Open) No. 68449/1986]. This liquid crystal compound has a very large spontaneous polarization of 180 nC/cm.sup.2 but, being an aliphatic chloro compound, has poor chemical stability. Hence, there was synthesized 4'-octylcarbonyloxy-4-[(s)-2-methoxy-(s)-3-methylpentyloxycarbonyl]bipheny l [Japanese Patent Application Kokai (Laid-Open) No. 228036/1987]. This compound has excellent chemical stability but has an insufficient spontaneous polarization of 17 nC/cm.sup.2."}
{"text": "1. Field of the Invention\nThe present invention relates to a process for the production of polyfluoroiodide.\nPolyfluoroiodides are very useful intermediates in producing compounds which are used as active components of a waterand oil-repellent, a mold-release agent, a finish and the like, through telomerization or epoxidation of an olefinic compound with the use of a peroxide.\n2. Description of the Related Art\nIn order to iodinate a fluorine-containing organic acid, for example, a process comprising forming a silver salt of said organic acid with for example a silver oxide and subsequently reacting the salt with iodine (U.S. Pat. No. 2,170,181, the Hunsdiecker reaction) and a process comprising reacting an alkali metal salt of the organic acid with iodine (G.B. Patent Application No. 86134) are proposed. However, in any process, there are some defects, for example, a reacting agent is expensive, building an apparatus in a commercial scale is difficult, or a yield is not satisfactory."}
{"text": "1. Field of the Invention\nThe present invention relates to a radiation imaging and therapy apparatus for obtaining a radiation image of a breast and having a therapy function of applying radiation to an affected part in the breast.\n2. Description of a Related Art\nRecent years, an image-guided radiation therapy (IGRT) apparatus combining a tomographic imaging apparatus and a radiation therapy apparatus has been developed. Thereby, a radiation application method can be planned based on a three-dimensional (3D) image near an affected part, and radiation therapy for applying radiation can be performed according to the plan.\nFor example, Japanese Patent Application Publication JP-A-9-192245 discloses a radiation therapy system as shown in FIG. 12. The radiation therapy system includes a bed movably supported in a direction in parallel with the reference surface, an imaging apparatus (X-ray CT apparatus) that can image a part for therapy of an object to be inspected on the bed, a positioning apparatus for positioning the part for therapy from the image obtained by the imaging apparatus, and a radiation therapy apparatus for performing therapy by applying radiation to the positioned part for therapy. The imaging apparatus, the positioning apparatus, and the radiation therapy apparatus are accommodated together in one room, and a patient as the object undergoes therapy while staying on the bed because the bed sequentially moves to the locations of the respective apparatuses.\nHowever, the image-guided radiation therapy apparatus can be used for all body organs as targets of examinations, and there is a problem that a contrast agent is necessary for an examination of breast cancer because an X-ray having high penetrating power generated by an X-ray tube voltage equal to or more than 100 kV for image formation is used. Further, since image-guided radiation therapy apparatus is used for a whole body as a target of examination, upsizing of the entire apparatus is unavoidable.\nOn the other hand, for the main purpose of detections and diagnoses of breast cancer, a radiation tomographic imaging apparatus for breast (mammography CT apparatus) that can obtain three-dimensional (3D) images has been developed. As an example of a mammography CT apparatus, U.S. Pat. No. 6,480,565 B1 discloses a device for producing a three-dimensional tomographic mammography image of a breast of a patient. The device disclosed in U.S. Pat. No. 6,480,565 B1 comprises: a gantry frame; at least one motor for moving the gantry frame to form a data acquisition geometry; a source of radiation attached to the gantry frame to move with the gantry frame; a flat panel detector attached to the gantry frame to move with the gantry frame, the flat panel detector being disposed in a path of the radiation; and a support on which the patient rests while the mammography projection images are taken, the support supporting the patient such that the breast is disposed between the source of radiation and the flat panel detector; the support comprising a table on which the patient lies while the mammography projection images are taken; wherein the at least one motor moves the gantry frame so that the flat panel detector takes a volume scan of the breast; and wherein the table has two breast holes for both of the patient's breasts. For imaging a tumor mass or calcification in the breast, an X-ray with high energy has a small ratio of transmission amount and the contrast of transmission image is weak, and therefore, an X-ray beam with low energy of about 25 kV to 50 kV is used.\nWhen breast cancer detection is performed by using the device disclosed in U.S. Pat. No. 6,480,565 B 1, the patient rests face down on the table to allow two breasts to extend through the two breast holes, and one breast descends between the radiation source and the flat panel detector. The radiation source and the flat panel detector perform radiation imaging at each predetermined angular position while rotating together around a rotational axis, and thereby, the flat panel detector obtains radiation images of the breast in plural directions. The image signals obtained by the flat panel detector are transmitted to an image reconstruction and processing module. The image reconstruction and processing module three-dimensionally reconstructs the obtained radiation images to create a radiation tomographic image of the breast. Since the radiation tomographic image of the breast is a 3D image, a focus part at the rear side of the tissues such as mammary gland can be detected.\nAccordingly, in the image-guided radiation therapy apparatus for performing radiation therapy for breasts, it is conceivable that the mammography CT apparatus is used in place of the imaging apparatus that can obtain radiation projection images in plural directions. However, when the mammography CT apparatus and the radiation application therapy apparatus are combined, it is difficult to reproduce the affected part location in the radiation application therapy apparatus for accurate therapy of the affected part in the breast, which has been recognized by using the mammography CT apparatus."}
{"text": "Plastics and/or rubbers and composites typically require the use of additives and/or fillers to prepare viable commercially useful materials. Additives and/or fillers are essential for their processing and have a direct impact on the cost-performance of the final product. Specific and desired aesthetic and/or color can also be targeted through the use of these materials.\nFillers are particulate additives used for a variety of purposes to modify the performance of plastics. Carbon black is widely used in considerable quantities as reinforcing filler in the rubber industry. However, when it is used in, for example commodity plastics, the quantities based on total weight percentage do not normally exceed 4% due to a loss in the mechanical properties of the doped plastic. Carbon black in plastics, for example for high density polyethylene piping or polyolefins for automotive applications does not typically exceed these percentages. In automotive parts' applications, carbon black is used as pigment and as an ultraviolet absorber in a variety of plastics including, but not limited to polypropylene (PP); polyethylene (PE); thermoplastic polyolefin (TPO) (rubberized polyolefin); or their blend; and poly (vinyl chloride) (PVC).\nWorld-wide production of commodity carbon black in 2011 was 14.26 million tones and is expected to grow 5% from 2013-2018. Carbon black is derived from non-renewable resources such as natural gas or petroleum-derived heavy oils through a chemical-thermal conversion. The soot formed in this process is recovered and processed into carbon black. The most common process today is the “furnace black process”. In this process each kilogram of carbon black requires approximately two and a half kilograms of oil. The fossil-fuel carbon black industry is a major contributor to carbon dioxide emissions, which has been implicated in global warming.\nAlthough elemental carbon in carbon black occurs at nanometer scale, aggregate states may occur at the micron scale. In addition, carbon black is produced as fine dust which represents a serious environmental and health risk and hazard if manipulated in this form. A common solution to this problem is to pelletize the carbon black in the presence of binding agents into agglomerates or ‘pearls’ or to produce master batches of carbon black prior to final dosage in plastics. However, those familiar with the technology of using “pearls” of carbon black master batch pellets are aware of the challenges that the plastic industry faces to disperse these carbon blacks during the final dosage for extrusion or injection molded pieces.\nThe ASTM standard D3350 Section 2 points out that the concentration of carbon black “in carbon-doped” plastics should be higher than 2% by weight in order to provide an adequate UV protection. However, as previously stated, it is known by those versed in this art that the use of carbon black should not exceed 4 wt % (more preferably 3% and even more preferably 2.5 wt %) of the total composition in the case of plastics such as polyolefins, PVC, or similar thermoplastics used for pipes and/or automotive applications. When carbon black is added in excess, it generally reduces the general mechanical properties of the plastics."}
{"text": "1. Field of the Invention\nThe invention relates to computer systems, and more particularly relates to computer systems having a voice-control function.\n2. Description of the Related Art\nElectronic systems normally have power management, and there are lots of standards of power management, such as the Advanced Configuration and Power Interface (ACPI) specification. The ACPI is a power-management open standard developed by Intel, Microsoft and Toshiba, and is suitable for all classes of computer systems including desktop computers, portable computers, workstations, and server machines. The purpose of the ACPI specification is to efficiently distribute power supply to each component in the computer system. It should be noted that ACPI is a power management interface shared between the software (operating system, OS) and hardware, and power management herein is OS-directed instead of BIOS-directed, which is more beneficial for manufacturers when integrating their standards with each other.\nAlso, with most computer systems, a user needs to operate them by direct touch. However, it is difficult for a user with disabilities to operate them. The technique of voice recognition has been under development for long time, and today the voice recognition rate is quite high. Therefore, using voice recognition to improve the operability of computer systems for users has become an important issue."}
{"text": "Devices within a distributed computing environment must typically exchange large amounts of configuration data in order to identify and access various distributed resources. For example, a central server within a client-server computing model may be required to distribute a large amount of configuration data to each client within its distributed computing environment in order to enable these clients to identify and access the distributed resources they require to perform various tasks. Similarly, a peer within a peer-to-peer computing model may request a large amount of configuration data from one or more of its peers in order to identify and access the distributed resources the peer requires to perform various tasks.\nA distributed computing system's configuration data is typically represented in the form of a tree that identifies the hierarchical relationship of each object within the distributed computing system. Nodes within this tree may represent objects (such as a cluster, system, group, application, resource, or attribute of the same) within the distributed computing system connected in parent-child relationships. Devices within the distributed computing system may be interested in either the entire tree (e.g., peers within a peer-to-peer model or a server within a client-server model may be interested in the entire tree) or a subset of the tree (e.g., clients within a client-server model may only be interested in a subset of the tree).\nUnfortunately, due to the typical size of this configuration data, an inordinate amount of network bandwidth may be required to exchange even a subset of this configuration data among devices within the distributed computing system, resulting in potential delays. For example, the amount of time required for a client to recover from a failure may be dramatically increased due to the time required to request, receive, and process a snapshot of (even a subset of) a central server's configuration data. As such, the instant disclosure identifies a need for systems and methods for efficiently synchronizing configuration data within distributed computing systems."}
{"text": "There are a variety of imaging techniques that have been used to diagnose disease in humans. One of the first techniques employed was X-rays. In X-rays, the images of the patients' body reflect the different densities of body structures. To improve the diagnostic utility of this imaging technique, contrast agents are employed to increase the density difference between various structures, such as between the gastrointestinal tract and its surrounding tissues. Barium sulfate and iodinated contrast media are the two primary pharmaceutical oral contrast agents used for X-ray gastrointestinal studies to visualize the esophagus, stomach, small intestine, and large intestine. Likewise, these same contrast agents are used for X-ray-based CT images to improve visualization and distension of the gastrointestinal tract and to provide improved contrast between the gastrointestinal tract and the structures adjacent to it, such as the blood vessels and lymph nodes. Such gastrointestinal oral contrast agents increase the density inside the esophagus, stomach, small intestine, and large intestine, and allow differentiation of the gastrointestinal system from the surrounding structures. Such gastrointestinal oral contrast agents also distend the bowel lumen allowing better evaluation of bowel wall structures, including abnormalities such as wall thickening or luminal obstruction.\nMagnetic resonance imaging is another imaging technique; however, unlike X-rays and CT, MR does not utilize ionizing radiation. MR employs a magnetic field, radiofrequency energy and magnetic field gradients to make images of the body. The contrast or signal intensity differences between tissues mainly reflect the T1 and T2 relaxation values and the proton density of the tissues. Like CT, MR can make cross-sectional images of the body and it is desirable for many MR scans that a gastrointestinal contrast agent be administered to allow differentiation of the gastrointestinal system from the surrounding structures. Such gastrointestinal oral contrast agents also distend the bowel lumen allowing better evaluation of bowel wall structures, including abnormalities such as wall thickening or luminal obstruction.\nOpacification of the gastrointestinal tract is a routine part of patient preparation prior to CT examination of the abdomen and pelvis because proper filling and distention of the bowel lumen with contrast medium facilitates identification of the normal and abnormal anatomy. For this reason, oral contrast agents are administered to the vast majority of patients undergoing CT scans of the abdomen and pelvis. Less commonly, oral contrast agents are administered to patients undergoing MR scans of the abdomen and pelvis.\nIn the early 1980s, when CT scan imaging was first gaining widespread use, the first oral agents used were the ionic iodinated contrast agents, comprised of the iodine-containing salts sodium diatrizoate, meglumine diatrizoate, or mixtures thereof. One example of the ionic iodinated oral contrast agents is Gastrografin® (Bracco Diagnostics, Princeton, N.J.), which is composed of the iodine-containing salts meglumine diatrizoate and sodium diatrizoate. Pharmaceutical agents like Gastrografin® were already approved for oral use and were widely used in full strength concentration for diagnostic gastrointestinal radiology imaging. These products, as well as the newer non-ionic iodinated agent approved for oral administration in the United States (e.g., iohexol, trade name Omnipaque®, GE Healthcare, Princeton, N.J.), were too concentrated in the form supplied by the manufacturers to be used for gastrointestinal opacification during CT imaging. These products were, therefore, diluted with water or other beverages in order to be of the proper concentration for optimal CT scan imaging.\nSimilarly, the gadolinium-based oral contrast agents, such as gadopentetate dimeglumine, are too concentrated in the form supplied by the manufacturers to be used for gastrointestinal opacification during MR imaging. These products must be diluted with water or other beverage in order to be of the proper concentration for optimal gastrointestinal opacification during MR imaging.\nIn addition to dilution for imaging, oral contrast agents are mixed with beverages to improve palatability and taste. Due to the bitter and unpleasant taste of oral contrast agents, and the relatively large volume of liquid that must be consumed, many patients have difficulty drinking an oral contrast agent unless it is mixed with a beverage that masks the unpleasant taste. Soda, fruit juices, flavored water or powdered drink mixes such as Kool-Aid™ or Crystal Light™ have been mixed with oral contrast agents to improve palatability and taste. In some cases, the beverage is poured, or in the case of powdered drink mixes combined with water, in a pitcher or similar container. The oral contrast agent is then added to the beverage in the pitcher and dispensed in a cup to the patient. The pitcher may contain the diluted oral contrast for a single patient or for multiple patients.\nUsing a pitcher and dispensing the oral contrast agent in a cup results in difficulty in monitoring and ensuring that the patient consumes the desired quantity of oral contrast agent over the time period established, particularly if the pitcher is used for oral contrast agent for more than one patient. It is also difficult to provide labeling or means to write on the pitcher the dose or type of oral contrast used, the patient name and medical identification number, and other data that may be of medical or regulatory significance.\nIn other cases, the dilute oral contrast agent is mixed and dispensed in a series of cups Because typically 32 ounces or more of the dilute oral contrast agent must be consumed by the patient, multiple cups are required. The use of multiple cups of oral contrast agent results in difficulty in monitoring the rate of consumption. In addition, labeling of the cups is typically minimal, and in some cases no labeling is used at all.\nAccordingly, it would be advantageous to have a container or other device for administration of oral contrast agent that allowed the user to obtain the proper concentration and dilution of the oral contrast agent, monitor consumption of the oral contrast agent to ensure that proper quantities are consumed within the time period established, and/or provide easy labeling to track information such as the patient name, oral contrast agent used, concentration of oral contrast agent in the container, date and time of preparation, and quantities to be consumed for specific time periods."}
{"text": "This disclosure relates to systems and methods for controlling systems and equipment that provide cooling for areas that contain electronic equipment, such as computer server rooms and server racks in computer data centers."}
{"text": "The ability to test the tensile strength and other mechanical properties of single fibers is of great importance for researchers, engineers, and manufacturers. Single fiber testing can be used to isolate trends and modes of failure for materials, such as for fundamental materials investigations or product-line quality control testing. During development of new fiber materials, the new fibers are often available in very small quantities due to complex processing or exotic materials. In these cases, multi-filament yarn-level testing is not practical. Instead, it is desired to generate mechanical properties of individual filaments.\nDevelopments in capacitive load cells, electromagnetic load cells, and stepper motor technology have led to the commercialization of compact load frames with extremely sensitive force (for example, nanoNewton (nN) sensitivity) and displacement (for example, nanometer (nm) sensitivity) measurement capabilities. These miniature load frames are ideally suited for testing single fibers.\nStandard testing procedures for single fibers[1] call for directly gripping single fibers in clamping action grips, or gluing the fibers onto cardboard tabs that are then gripped in clamping action grips. The first technique, direct clamping, concentrates stresses at the grip point and often leads to failure at the fiber grips[2]. Grip failure is undesirable during fiber characterization, and these results are generally not accurate, and, thus, are not reportable. The second technique, bonding fibers onto cardboard tabs, relies on strong bonding between the fiber and a dollop of adhesive. For fibers that are difficult to bond to, such as aramid, polyethylene, or ultrahigh molecular weight polyethylene, it is common to observe fiber slippage from the adhesive. These results are also not accurate and, thus, are not reportable, and generally prevent loading the fiber to failure. In addition, the time and complication associated with utilizing curable adhesives adds considerable inconvenience, especially considering that many test repetitions are required to establish statistically significant data. Accordingly, there is a need for a method, a device and/or an assembly for clamping or otherwise securing a fiber in a tensile testing device.\nFor testing yarns, it is common to use capstan grips (also called roller grips) or contour/horn grips (also called half-capstan grips)[2, 3]. In these designs, from the primary loading gage section of the yarn the yarn is passed over a series of “snubbing” surfaces which gradually reduce the tensile load on the yarn due to frictional interactions of the yarn with the snubbing surfaces. The snubbing surfaces are curved gently to induce normal forces on the yarn without creating stress concentrations that can fail the yarn. After passing over these snubbing surfaces, the yarn is gripped mechanically, typically using mechanical or pneumatic clamping action. Because the tensile loading on the yarn has been gradually reduced, the total stress on the yarn at the grip point is significantly less likely to lead to failure at the grip.\nThe present invention improves upon previous single fiber techniques by providing a single fiber clamping device that can reliably grip single fibers and decrease slippage and/or failure at the grips. In addition, no adhesive is required, so that gripping is fast and uncomplicated. Furthermore, the device uses magnetic clamping action to grip the fibers, so that the total mass and volume of the clamping device are minimized. Maintaining low clamping device mass is desirable for operation in sensitive load frames, since the total load limits (including clamp weight) are typically very low for these miniature load frames (for example, 1 N maximum force capacity). For these applications, traditional clamping approaches such as mechanical action or pneumatics are undesirable because they are either impractical or cumbersome to implement."}
{"text": "1. Field of the Invention\nThis invention relates to the field of cords or cables. More specifically, the present invention comprises a cord construction, a cord termination, and an adjustable clasp which can be used to attach a cord to another device.\n2. Description of the Related Art\nVarious fastening mechanisms are known for attaching cords or cables to devices (Throughout this disclosure, the terms “cord” and “cable” are considered synonymous). Fastening mechanisms are particularly important when the cord supports a tensile load. Most fastening mechanisms are not adjustable and do not allow the user to easily vary the length of the cord. These fastening mechanisms are typically semi-permanent. Often the user can shorten the length of the cord by cutting the cord next to the fastener and reattaching the same fastener or a different fastener to the shortened cord. It is generally inadvisable if not impossible to lengthen the cord because conventional fastening mechanisms damage the portion of the cord at the location where the fastener is attached.\nOne example of applications where it is useful to have an adjustable fastening means is control cable applications. For purposes of illustration, a bicycle control cable, such as a control cable used for braking or shifting, will be considered. A complete control cable consists of the tensile carrying cable itself, a first terminating device on one end, and a second terminating device on the other end. Because the length of such cables often need to be adjusted, at least one of the two terminating devices will normally include an adjustment feature.\nBicycle control cables are typically composed of an inner steel cable which is wrapped in a protective and supportive outer sheath. Various fastening and actuating devices are customarily attached to the opposite ends of the inner steel cable. These can include combinations of crimps, ferrules, donuts, caps and other hardware. Such fastening devices are used to connect the inner steel cable to a braking handle (or shifting lever) on one end and to the braking mechanism (or shifting mechanism) on the other end. Generally, some piece of hardware is crimped on to both ends of the inner steel cable so that the cable can be fastened to these other devices.\nCrimping is a semi-permanent method of attaching hardware to the inner steel cord. To remove the hardware, the steel inner cord must generally be cut beneath the crimp. If the user attaches the hardware too short to reach the device it is designed to actuate or if the user moves the actuating mechanism further from the device that is actuated, the user will often have to replace the control cable. Accordingly it would be desirable to have an adjustable fastening mechanism that can be attached without damaging the cord.\nIn recent years the traditional steel cable has been replaced by advanced synthetic fibers. These fibers, when compared to steel, are much smaller and slicker. Crimping does not work well for synthetics, since the crimping force required to grip such a slick surface will often damage the fibers. Synthetic fibers also have very little compressive stiffness. This fact means that a bundle of such fibers will tend to be loosely organized. It will not maintain a circular cross section when passed around a bend but will flatten instead—producing uneven load distribution across the cable's cross section. An integrated cable, termination, and clasp system should ideally take into consideration the new design constraints imposed by the growing use of synthetic fibers."}
{"text": "1. Field of the Invention\nThe invention relates to a communication network and the handling of service connections in the communication network including an Internet access system and a switching center.\n2. Description of Related Art\nA strong growth of the traffic volume in the access area to services that are offered by other networks, for example the Internet, is connected with the strong growth of Internet services and other special services. As a result, there is a substantial increase in the demands made of the exchanges forming the access area to these networks.\nThe situation is additionally intensified in that current exchanges are not dimensioned for the absolute call lengths and distributions of call lengths that mainly occur in the Internet traffic, to the entry nodes of the Internet service providers. Even when there is a slight increase in Internet traffic volume, additional exertions of the Internet access service operators are needed for expanding the capacity of their switching centers in order to maintain the actual voice traffic with the required quality features."}
{"text": "There is currently an ongoing drive toward the downscaling of device dimensions in virtually all aspects of electronic device manufacture. Smaller electronic devices tend to be more popular than larger, more bulky devices when both devices have substantially equivalent capabilities. Accordingly, being able to fabricate smaller components would clearly tend to facilitate the production of smaller devices that incorporate those components. However, many modern electronic devices require electronic circuitry to perform both actuation functions (e.g., switching devices) and data processing or other decision making functions. The use of low voltage complementary metal-oxide-semiconductor (CMOS) technologies for these dual functions may not always be practical. Thus, high voltage (or high power) devices have also been developed to handle many applications where low voltage operation is not practical.\nThe electrostatic discharge (ESD) performance of typical high voltage devices often depends on the total width and surface or lateral rules of the corresponding devices. Thus, ESD performance may typically be more critical for smaller devices. High voltage devices typically have characteristics that include a low on-state resistance (Rdson), a high breakdown voltage and a low holding voltage. The low on-state resistance may tend to make an ESD current more likely to concentrate on the surface or the drain edge of a device during an ESD event. High current and high electric fields may cause the physical destruction at a surface junction region of such a device. Based on the typical requirement for a low on-state resistance, the surface or lateral rules likely cannot be increased. Thus, ESD protection may be a challenge.\nThe high breakdown voltage characteristic of high voltage devices typically means that the breakdown voltage is higher than the operating voltage, and the trigger voltage (Vt1) is higher than the breakdown voltage. Accordingly, during an ESD event, the internal circuitry of the high voltage device may be at risk of damage before the high voltage device turns on for ESD protection. The low holding voltage characteristic of high voltage devices also leaves open the possibility that unwanted noise associated with a power-on peak voltage or a surge voltage may be triggered or that a latch-up may occur during normal operation. High voltage devices may also experience the field plate effect due to the fact that electric field distribution may be sensitive to routing so that ESD current may be likely to concentrate at the surface or drain edge during an ESD event.\nTo improve high voltage device performance with respect to ESD events, one technique that has been implemented involves the additional use of masks and other processes to create a larger sized diode within bipolar junction transistor (BJT) components and/or increasing the surface or lateral rules for MOS transistors. Silicone controlled rectifiers (SCRs) have also been developed to protect circuitry during ESD events. However, while the low holding voltage of SCRs means they may perform well during ESD events, this characteristic also increases the occurrence of latch-up during normal operation.\nMotor driver circuits may be particularly troublesome to protect from ESD events using current solutions. This is because when a motor is switched off, it may continue spinning for some time, thus acting as an inductor which feeds back a high negative voltage. If the motor driver circuitry were to include a PMOS, the parasitic forward bias diode of the PMOS may be turned on by this negative feedback voltage, potentially causing latch-up issues and/or other irregular circuit operation.\nAccordingly, it may be desirable to develop an improved structure for providing ESD protection and, in particular, for providing bi-directional ESD protection."}
{"text": "The heart is surrounded by the pericardial sac, a loose layer of tissue, which may be visualized and retracted from the heart in an open chest. Retracting the pericardial sac from the heart creates the pericardial space, a space between the heart and the pericardial sac. U.S. Pat. No. 6,162,195 issued to Igo, et al. describes the pericardial sac as surrounding the heart like a glove enfolds a hand, and the pericardial space as naturally fluid-filled. The normal pericardium functions to prevent dilatation of the chambers of the heart, lubricates the surfaces of the heart, and maintains the heart in a fixed geometric position. The normal pericardial space is small in volume and the fluid film within it is too thin to functionally separate the heart from the pericardium. See Shabetai R: Pericardial and cardiac pressure, Circulation 77:1, 1988\nMany cardiac surgical procedures invade the pericardial sac leaving it torn with no detrimental effect. Various approaches may be used to access the pericardial sac and the outside of the heart including a median sternotomy, a thoracotomy, a thoracoscope, a sub-xiphoid route and a transvenous route. The median sternotomy is a chest surgical technique using an incision of the skin and sternum between the xiphoid process and the suprasternal notch. The thoracotomy is a chest surgical technique using an incision of the chest wall and may include removal of a portion of a rib. The thoracoscope approach generally involves three puncture wounds to the chest, placement of ports for access with light, camera and instruments. The sub-xiphoid route uses a puncture inferior to the xiphoid process and the xiphoid process may or may not be removed. Transvenous approaches to the outside of the heart generally involve venous access on the right side of the heart and a puncture through the right atrium, right atrial appendage or superior vena cava.\nAccess to the outside of the heart may be useful for revising the circulation of the heart as in coronary artery bypass grafting (CABG), the delivery of pharmaceutical agents, the revision of the heart structure in support of cardiac valvular function, the placement of electrodes for pacing, sensing, monitoring, cardioverting or defibrillating, the ablation of tissue for the prevention of arrhythmias, or other structural modification. Implantable stimulation of the heart has been in clinical use since the early 1960's when electrodes were implanted on the outside of the heart, the epicardium. Insulated wires connecting the electrodes to an implantable pulse generator called “leads” were implanted by a surgical procedure.\nApproaches that involve an intact chest such as the sub-xiphoid do not provide space for accessing and retracting the pericardial sac. To access the outside of the heart without unintentional perforation, incision or damage to the heart requires retraction of the pericardial sac. U.S. Pat. No. 6,315,774 issued to Daniel et al. discloses a thoracoscopic approach to access the space around the heart by the use of a pair of graspers and scissors, then grabbing and opening the pericardial sac by making a stem to stern type of incision. The pericardial sac is pulled away from the heart and may be suspended.\nU.S. Pat. No. 7,597,698 issued to Chin describes creating an opening through a pericardial reflection, a fold in the pericardium but is quick to point out that dissection is hazardous because of important large blood vessels in the vicinity of the reflection. Grasping and forming a hole in the superior vena cava, for example, would be disastrous. U.S. Pat. Nos. 4,181,123 and 4,319,562 to Crosby, and 5,033,477 to Chin et al. disclose methods for placing electrodes in contact with the heart muscles from within the pericardial space without the need for thoracotomy via a sub-xiphoid route which involves penetrating the chest wall below the xiphoid process. However, Waxman, in U.S. Pat. No. 5,968,010, describes the sub-xiphoid route as so small that it is difficult to penetrate the sac without also puncturing, and thereby, damaging the heart itself. U.S. Pat. No. 4,991,578 issued to Cohen discloses distending the pericardium from the heart by injecting a small volume of fluid into the pericardium, puncturing the pericardium with a needle, passing a guide wire through the needle into the pericardium and then removing the needle. Adding fluid between the beating heart and the pericardial sac decreases the volume to which the heart may fill. Caution must be taken to avoid the extreme of this situation called tamponade in which the output of the heart is so restricted that systemic circulation is impaired.\nU.S. Pat. No. 6,237,605 issued to Vaska, et al. describes devices and methods that may be utilized through a small access port in the chest, preferably through a subxiphoid penetration, and positioned within the pericardium and around the pulmonary veins without cutting or puncturing the pericardial reflections.\nCryogenic techniques applied to the heart demonstrate reversible changes to the myocardium. U.S. Pat. No. 5,733,280 issued to Avitall describes the cooling of cardiac cells and rewarming the tissue resulting in total recovery of the tissue without damage. U.S. Pat. No. 5,147,355 issued to Friedman, et al. discloses a catheter having a fluid flow passage for directing a flow of cryogenic fluid to the tip of the catheter.\nStabilization with devices applied to a beating heart is recognized in U.S. Pat. No. 6,960,205 issued to Jahns, et al., incorporated herein in its entirety by reference, describing devices that use vacuum or suction force to hold tissue. Apparatus for visualization and access is described in U.S. Pub. No. 2007/0293724 to SAADAT, et al. wherein tissue can be engaged using a vacuum or a cryo-probe. Creation of an ice-ball at the catheter tip that stabilizes the tip relative to a tissue is described in U.S. Pub. No. 2007/0116921 by Sherman, et al.\nThe use of suction for accessing an anatomical space of the body, and particularly for penetrating the epicardium to access pericardial space and the epicardial surface of the heart is described in U.S. Pat. No. 6,890,295 issued to Michels et al., incorporated herein in its entirety by reference. The use of suction requires a relatively large diameter apparatus to achieve sufficient retention force applied to the pericardial surface as compared to the diameter of cardiac rhythm leads.\nU.S. Pat. No. 3,737,579 issued to Bolduc, incorporated herein in its entirety by reference, describes electrical leads for myocardial implantation with a rigid helix on the distal end of the lead serving as an electrode screwed into body tissue. U.S. Pat. No. 4,142,530 issued to Wittkampf, describes the screw-in type of electrode as requiring sufficient room to approach the heart wall from a direction that is more of less normal or perpendicular to the surface to allow the helix to be screwed into the heart muscle.\nFor applications where it is desired to access the epicardial surface of the heart, what is needed is a procedure that does not have significant morbidity, does not compromise the patient during the procedure, only requires a small wound and carries little risk of perforating the heart or other major blood vessels."}
{"text": "Substrates with electronically active components distributed over the extent of a substrate are used in a variety of electronic systems, for example, flat-panel imaging devices such as flat-panel liquid crystal or organic light emitting diode (OLED) display devices and in flat-panel solar cells. A variety of methods may be used to distribute electronically active circuits over substrates, including forming the electronically active circuits on a substrate and forming the components on separate substrates and placing them on a substrate. In the latter case, a variety of assembly technologies for device packaging can be used.\nThe electronically active components are typically formed on a substrate by sputtering a layer of inorganic semiconductor material or by spin-coating organic material over the entire substrate. Inorganic semiconductor materials can be processed to improve their electronic characteristics, for example amorphous silicon can be treated to form low-temperature or high-temperature poly-crystalline silicon. In other process methods, microcrystalline semiconductor layers can be formed by using an underlying seeding layer. These methods typically improve the electron mobility of the semiconductor layer. The substrate and layer of semiconductor material can be photo-lithographically processed to define electronically active components, such as transistors. Such transistors are known as thin-film transistors (TFTs) since they are formed in a thin layer of semiconductor material, typically silicon. Transistors may also be formed in thin layers of organic materials. In these devices, the substrate is often made of glass, for example Corning Eagle® or Jade® glass designed for display applications.\nThe above techniques have some limitations. Despite processing methods used to improve the performance of thin-film transistors, such transistors may provide performance that is lower than the performance of other integrated circuits formed in mono-crystalline semiconductor material. Semiconductor material and active components can be provided only on portions of the substrate, leading to wasted material and processing costs. The choice of substrate materials can also be limited by the processing steps necessary to process the semiconductor material and the photo-lithographic steps used to pattern the active components. For example, plastic substrates have a limited chemical and heat tolerance and do not readily survive photo-lithographic processing. Furthermore, the manufacturing equipment used to process large substrates with thin-film circuitry is relatively expensive. Other substrate materials that may be used include quartz, for example, for integrated circuits using silicon-on-insulator structures as described in U.S. Patent Application Publication No. 2010/0289115 and U.S. Patent Application Publication No. 2010/0123134. However, such substrate materials can be more expensive or difficult to process.\nOther methods used for distributing electronically functional components over a substrate in the circuit board assembly industry include, for example, pick-and-place technologies for integrated circuits provided in a variety of packages, for example, pin-grid arrays, ball-grid arrays, and flip-chips. However, these techniques may be limited in the size of the integrated circuits that can be placed.\nIn further manufacturing techniques, a mono-crystalline semiconductor wafer is employed as the substrate. While this approach can provide substrates with the same performance as integrated circuits, the size of such substrates may be limited, for example, to a 12-inch diameter circle, and the wafers are relatively expensive compared to other substrate materials such as glass, polymer, or quartz.\nIn yet another approach, thin layers of semiconductor are bonded to a substrate and then processed. Such a method is known as semiconductor-on-glass or silicon-on-glass (SOG) and is described, for example, in U.S. Pat. No. 7,605,053, issued Oct. 20, 2009. If the semiconductor material is crystalline, high-performance thin-film circuits can be obtained. However, the bonding technique and the processing equipment for the substrates to form the thin-film active components on large substrates can be relatively expensive.\nPublication No. 11-142878 of the Patent Abstracts of Japan entitled Formation of Display Transistor Array Panel describes etching a substrate to remove it from a thin-film transistor array on which the TFT array was formed. TFT circuits formed on a first substrate can be transferred to a second substrate by adhering the first substrate and the TFTs to the surface of the second substrate and then etching away the first substrate, leaving the TFTs bonded to the second substrate. This method may require etching a significant quantity of material, and may risk damaging the exposed TFT array.\nOther methods of locating material on a substrate are described in U.S. Pat. No. 7,127,810. In this approach, a first substrate carries a thin-film object to be transferred to a second substrate. An adhesive is applied to the object to be transferred or to the second substrate in the desired location of the object. The substrates are aligned and brought into contact. A laser beam irradiates the object to abrade the transferring thin film so that the transferring thin film adheres to the second substrate. The first and second substrates are separated, peeling the film in the abraded areas from the first substrate and transferring it to the second substrate. In some embodiments, a plurality of objects is selectively transferred by employing a plurality of laser beams to abrade selected area. Objects to be transferred can include thin-film circuits.\nU.S. Pat. No. 6,969,624 describes a method of transferring a device from a first substrate onto a holding substrate by selectively irradiating an interface with an energy beam. The interface is located between a device for transfer and the first substrate and includes a material that generates ablation upon irradiation, thereby releasing the device from the substrate. For example, a light-emitting device (LED) is made of a nitride semiconductor on a sapphire substrate. The energy beam is directed to the interface between the sapphire substrate and the nitride semiconductor releasing the LED and allowing the LED to adhere to a holding substrate coated with an adhesive. The adhesive is then cured. These methods, however, may require the patterned deposition of adhesive on the object(s) or on the second substrate. Moreover, the laser beam that irradiates the object may need to be shaped to match the shape of the object, and the laser abrasion can damage the object to be transferred. Furthermore, the adhesive cure takes time, which may reduce the throughput of the manufacturing system.\nOther methods for transferring active components from one substrate to another is described in AMOLED Displays using Transfer-Printed Integrated Circuits published in the Proceedings of the 2009 Society for Information Display International Symposium Jun. 2-5, 2009, in San Antonio Tex., US, vol. 40, Book 2, ISSN 0009-0966X, paper 63.2 p. 947. In this approach, small integrated circuits are formed over a buried oxide layer on the process side of a crystalline wafer. The small integrated circuits, or chiplets, are released from the wafer by etching the buried oxide layer formed beneath the circuits. A PDMS stamp is pressed against the wafer and the process side of the chiplets is adhered to the stamp. The chiplets are pressed against a destination substrate or backplane coated with an adhesive and thereby adhered to the destination substrate. The adhesive is subsequently cured. In another example, U.S. Pat. No. 8,722,458 entitled Optical Systems Fabricated by Printing-Based Assembly teaches transferring light-emitting, light-sensing, or light-collecting semiconductor elements from a wafer substrate to a destination substrate or backplane.\nIn such methods, it is generally necessary to electrically connect the small integrated circuits or chiplets to electrically conductive elements such as backplane contact pads on the destination substrate. By applying electrical signals to conductors on the destination substrate the small integrated circuits are energized and made operational. The electrical connections between the small integrated circuits and the backplane contact pads are typically made by photolithographic processes in which a metal is evaporated or sputtered onto the small integrated circuits and the destination substrate to form a metal layer, the metal layer is coated with a photoresist that is exposed to a circuit connection pattern, and the metal layer and photoresist are developed by etching and washing to form the patterned electrical connections between the small integrated circuits and the connection pads on the destination substrate. Additional layers, such as interlayer dielectric insulators can also be required. This process is expensive and requires a number of manufacturing steps. Moreover, the topographical structure of the small integrated circuits over the destination substrate renders the electrical connections problematic. For example, it can be difficult to form a continuous conductor from the destination substrate to the small integrated circuit because of the differences in height over the surface between the small integrated circuits and the destination substrate.\nThere is a need, therefore, for robust structures and methods that enable the electrical interconnection of small integrated circuits, such as micro-transfer printed chiplets, to destination substrates."}
{"text": "Since the anatomic and physiologic structure of the anorectal section carries too specific features, treatment of the fistulas occurring in this area brings with it many difficulties. Taking into account that intersecting external sphincter and—puborectal—muscles will cause anal incontinence during the treatment, it is necessary to maintain the integrity of these muscles to a large extent. Anal fistulas can be examined under two categories, simple (can be treated with fistulotomy) and complex (when the fistulotomy contradicts/counteracts). Basic fistula can be treated easily by fistulotomy and—can be left—for secondary intention. In complex fistulas, set-on application or periodic surgeries are the methods used in order to remove the fistula trace.\nThis application is not only persistent but also hard for the patient and the treatment concludes with failure due to the high rate of incontinence risk. The treatment principle on fistulas, removing the dead tissues at the fistula trace, moving away the excreting mucous and opening the bloody tissue in order to stimulate the granulation tissue which helps filling the trace.\nAlthough this purpose is achievable with fistulotomy in basic fistulas, fistulotomy cannot be applied in complex fistulas as it will cause anal inefficiency. By gradually cutting the muscle tissue, which forms the sphincter, a method will be applied that will provide healing fibrous tissue.\nThe cut is executed while the reaction starts with cutting the muscle mass which develops granulation tissue, merging the execrated muscle tissues with fibrous structure and the expectation of filling the fistula trace. While this is not always successful, it causes loss of continence in a certain ratio due to the sphincter dysfunction which fuses with the fibrous tissue.\nHence, there is a need for a device to cure such fistulas."}
{"text": "In many wall designs, such as retaining walls or counterfort retaining walls or flood walls, the site parameters dictate the final design of the wall. In many cases there is also a property line that abuts the wall and moving the wall as close as possible to the property line is desirable. The retaining wall of the present invention gives the designer the ability of doing so. In addition, the embodiments of the retaining wall disclosed can be designed to have less concrete and backfill and thus may be more economical to construct."}
{"text": "1. Field of the Invention\nThis invention relates to a display device that has a plurality of pixels disposed in a matrix configuration, a plurality of gate lines extending in a row direction, and a driving circuit sequentially supplying a gate scanning signal to each of the gate lines, especially to a display device with improved image quality.\n2. Description of the Related Art\nFIG. 10 is a diagram of a conventional liquid crystal display device. A liquid crystal panel 100 has a plurality of pixels formed in a matrix configuration of n-rows and m-columns. Each of the pixels has a pixel selecting thin film transistor 10, a liquid crystal LC, and a storage capacitor Csc. The thin film transistor will be referred to as TFT hereinafter.\nA gate line 20 extending in a row direction is connected to the gate of the TFT 10, and a data line 22 extending in a column direction is connected to the drain of the TFT 10. A gate scanning signal is sequentially supplied from a vertical driving circuit (V drive circuit) 130 to the gate line 20 of each row, and the TFT 10 is selected accordingly. Also, a video signal is applied to the liquid crystal LC through the TFT 10 based on a drain scanning signal supplied from a horizontal driving circuit (H drive circuit) 140.\nJapanese Patent Application Publication No. Hei 10-115839 discloses a liquid crystal display device with the above configuration.\nHowever, a gate insulating layer of the TFT in the output portion of the gate scanning signal of the vertical driving circuit 130, sometimes suffers from a break down and leakage due to static charge induced during the TFT manufacturing process of the conventional liquid crystal display device. This kind of trouble will be explained by referring to FIG. 11A and FIG. 11B. FIG. 11A and FIG. 11B are enlarged view of part B encircled by the dotted line in FIG. 10. These figures show the pattern of the edge of the gate line 20 as well as the output portion of the vertical driving circuit 130. FIG. 11A is a plan view, and FIG. 11B is a cross-sectional view along with the X—X line in FIG. 11A.\nA dry-etching method is employed for processing the gate line 20, a storage capacitor line 21, and a gate wiring in the vertical driving circuit 130. Static charge is induced and stored in the gate line 20, the storage capacitor line 21, and the gate wiring during the dry-etching process. When the ion implantation of an N-type impurity such as arsenic or phosphorus into a P-Si layer, or the ion implantation of a P-type impurity such as boron in case of P-type channel TFT, is performed for forming a source region and a drain region using the gate line 20 as a mask, static charge is also stored in the gate line 20, the storage capacitor line 21, and the gate wiring due to the charge-up phenomena. It is especially easy for the gate line 20 and the storage capacitor line 21 to store static charge because they extend across the liquid crystal panel 100.\nThe stored charge is discharged from the edge of the gate line 20 to the gate wiring 13 of the TFT 1 located close to the gate line 20. The charges from the discharge will reach a gate electrode of an adjacent TFT 2 through the gate wiring 13, causing a dielectric break down and leakage at part A of the gate insulating layer 12 of the TFT 2-1. In the figure, the reference numeral 40 indicates a transparent insulating substrate, the reference numeral 11 denotes an active layer (poly-silicon layer) disposed on the transparent insulating substrate 40, and the reference numeral 14 indicates an aluminum wiring layer. FIG. 11B is a cross-sectional view of the device before the aluminum wiring layer 14 is disposed."}
{"text": "A self-optimizing network (SON) refers to a network performing self optimization in a certain manner/technique with no need of participation of an operation, administration and maintenance (OAM) entity.\nMobility robustness optimization (MRO) is SON means, with a main object being to detect and correct a connection failure of radio resource control (RRC) during a user moves resulted from improper setting of a network parameter.\nCurrently, connection failure referred in the MRO includes the following three types.\nFirst Type: Too-Late Handover\nThe too-late handover refers to that as the handover is triggered too late, a connection failure has occurred when a handover command has not been successfully transmitted to user equipment (UE). After the failure occurs, the UE will initiate connection reestablishment to a reestablished cell (if the network side has finished handover preparation), or will initiate connection reestablishment to other non-source cells (if the network side has not performed handover preparation).\nFIG. 1 is a schematic diagram of too-late handover. As shown in FIG. 1, as a parameter (such as a signal difference between a target cell signal and a source cell signal) is set improperly, that is, the signal difference Δ′ between the target cell signal and the source cell signal is set to be relatively large, which is larger than a signal difference Δ between the target cell signal and the source cell signal of an ideal handover point, thereby resulting in too-late handover.\nSecond Type: Too-Early Handover\nThe too-early handover refers to that as the handover is triggered too early, a connection failure occurs resulted from a handover is performed by UE when a signal of the source cell is still strong and a signal of the target cell is still weak. The connection failure may occur shortly after the UE is handed over to the target cell. And after the failure occurs, the UE initiates connection reestablishment to the source cell.\nFIG. 2 is a schematic diagram of too-early handover. As shown in FIG. 2, as a parameter (such as a signal difference between the target cell signal and the source cell signal) is set improperly, that is, the signal difference Δ′ between the target cell signal and the source cell signal is set to be relatively small, which is smaller than a signal difference Δ between the target cell signal and the source cell signal of an ideal handover point, thereby resulting in too-early handover.\nThird Type: Handover to Wrong Cell\nThe handover to wrong cell refers to that, due to improper setting of a handover parameter, UE which should be handed over to a cell A is handed over to a cell B; and at this moment, as signal quality of the cell B is not good, a connection failure is resulted. The connection failure may occur shortly after the UE is handed over to the cell B. And after the failure occurs, the UE initiates connection reestablishment to target cell A.\nAfter the connection failure occurs when the UE is handing over, the UE will attempt to perform RRC connection reestablishment. If the connection reestablishment succeeds, the UE will keep the state of connection with the network side, and acquire services from the network side; and if the connection reestablishment fails, the UE will enter into an idle state, the connection with the network side will be interrupted, and the UE will possibly establish connection with the network side by initiating a new RRC connection reestablishment procedure.\nThe method for judging type of a handover failure in the MRO shall be described below regarding two cases where UE restores connection with the network side through an RRC reestablishment procedure or by initiating a new RRC connection establishment procedure.\nFirst: After the UE Initiates RRC Connection Reestablishment\n1. Too-Late Handover\nA radio link failure (RLF) occurs when the UE is in a source cell A, and the UE initiates a connection reestablishment attempt to a cell B, and transmits a connection reestablishment request (RRC) message to the cell B, the request message may include identifier information of the UE and an identifier of a cell where the UE is present when the RLF occurs, that is, an identifier of the source cell A. After receiving the request message, the cell B may transmit an RLF indication message to the source cell A where the RLF occurs according to information in the request message, the indication message including information contained in the RRC connection reestablishment request message; and after receiving the RLF indication message, the source cell A may judge type of the handover failure according to the indication message and saved context information of the UE; if the UE does not receive a handover command according to the context information of the UE, the type of the handover failure may be judged as too-late handover.\n2. Too-Early Handover\nThe UE is successfully handed over from the source cell A to the cell B, after a short period of time, an RLF occurs in the UE, and the UE initiates reestablishment connection to the cell A. During the procedure of the reestablishment connection, the UE transmits information, such as a user identifier and an identifier of the cell B where the UE is present when the link failure occurs to the cell A; the cell A transmits an RLF indication message to the cell B according to the obtained information; and if the cell B finds that a duration of the UE staying in the cell B is less than a parameter Tstore_UE_cntxt after receiving the message, it is judged that the UE turns back to the source cell A shortly after being connected to the cell B, and the cell B judges that too-early handover occurs in the UE.\n3. Handover to Wrong Cell\nThe UE is successfully handed over from the source cell A to the cell B, but a connection failure occurs after a short period of time. The UE initiates a reestablishment connection to a cell C, and at this moment, the cell C transmits an RLF indication message to the cell B, information contained in the message being similar to that as described above; and if the cell B finds that a duration of the UE connected to the cell B is less than a parameter Tstore_UE_cntxt after receiving the message, it is judged that the UE is connected to the cell C shortly after being connected to the cell B, and the cell B judges that handover to wrong cell occurs in the UE.\nSecond: After the UE Initiates a New RRC Connection Reestablishment Procedure\n1. Too-Late Handover\nAn RLF occurs when the UE is in the source cell A, and the UE initiates a connection reestablishment attempt to the cell B and the reestablishment fails. Then the UE initiates an RRC connection establishment procedure to the cell C (a cell selected via cell reselection, which may also be the cell A or B) via a cell reselection procedure, and reports an RLF report, the RLF report including an identifier of the cell where the UE is present when the RLF occurs, that is, an identifier of the source cell A, and sometimes further including a parameter timeConnFailure, which records a period of time experienced by the UE from receiving a handover command message for the last time to the occurrence of the link failure. The cell C transmits the RLF report reported by the UE via an RLF indication message to the source cell A where the RLF occurs. And if the source cell finds that the parameter timeConnFailure is not included in the RLF report or a value of the parameter timeConnFailure is greater than a preset value Tstore_UE_cntxt, the source cell A records this time of connection failure event as too-late handover.\n2. Too-Early Handover\nThe UE is successfully handed over from the source cell A to the cell B, after a short period of time, an RLF occurs in the UE, and the UE initiates a connection reestablishment attempt to the source cell A and the reestablishment fails. Then the UE initiates an RRC connection establishment procedure to the cell C (a cell selected via cell reselection, which may also be the cell A or B) via a cell reselection procedure, and reports an RLF report, the RLF report including an identifier of the cell B where the UE is present when the RLF occurs and the parameter timeConnFailure. The cell C transmits the RLF report to the cell B via an RLF indication message. And after the cell B obtains the RLF report, if the parameter timeConnFailure is less than Tstore_UE_cntxt and the reestablished cell of the UE is a cell triggering the latest handover before the connection fails, it is judged that this time of connection failure event is too-early handover.\n3. Handover to Wrong Cell\nThe UE is successfully handed over from the source cell A to the cell B, but an RLF occurs after a short period of time, and the UE initiates a connection reestablishment attempt to the cell C and the reestablishment fails. Then the UE initiates an RRC connection establishment procedure to a cell D via a cell reselection procedure, and reports an RLF report. The UE initiates reestablishment connection to the cell D, and at this moment, the cell D transmits an RLF indication message to the cell B, information contained in the message being similar to that as described above; and if the cell B, after receiving the message, finds that the parameter timeConnFailure is less than the parameter Tstore_UE_cntxt and the reestablished cell C of the UE is neither a source cell for the latest handover before the connection fails, nor a target cell, it is judged that this time of connection failure event is handover to wrong cell.\nHowever, in the implementation of the present disclosure, the inventors found that a defect exists in the prior art is: for a heterogeneous network deployed with a relatively small cell (such as a pico, and an RRH, etc.), it is possible that judgment error will occur, thereby resulting in that setting parameters cannot be adjusted in time, and the network performance is affected."}
{"text": "1. Field of the Invention\nThis invention relates to clamping devices.\nMore particularly, the present invention relates to devices normally employed as clamps on generally cylindrical pipes or bars.\nIn a further and more specific aspect, the present invention relates to retaining weights on weightlifting and barbell devices.\n2. The Prior Art\nDevices which are used as clamps on pipes, bars, shafts and the like are well known in the art. Specifically, devices for retaining weights on a bar for weightlifting applications are well known. Generally, a bar is provided onto each end of which disk shaped weights are mounted alone or in combination. Each disk shaped weight has a hole through its center, through which the bar is inserted. The weight disks are prevented from sliding off the bar by a collar which is placed on the bar between the outermost weight disk and the end of the bar. There are a large variety of collars, each preventing the weight disks from moving.\nWhen weightlifting devices of the barbell type are used, it is very important that the disk shaped weights are securely held on the bar. Considerable damage could occur if a weight disk were to become detached while the barbell was in use. Many of the prior art collars are secured to the bar with a screw extending through the collar. When the screw is tightened, it engages with the bar, locking the collar in place. While these clamps are generally effective, since the screw engages a very small area of the bar, it must be very tight to be adequately secured. This frequently makes the clamp difficult to remove, even when the screw has a lever to facilitate engaging and disengaging the collar. This may be a problem, since many times, a collar must be removed frequently so the weight configuration can be changed. This can occur because an exercise program requires different weights, or more than one person is using the barbell and each requires different weights.\nIt would be highly advantageous, therefore, to remedy the foregoing and other deficiencies inherent in the prior art.\nAccordingly, it is an object of the present invention to provide a quick releasing clamp.\nAnother object of the present invention is to provide a barbell clamp which securely engages the bar, but is easily removed."}
{"text": "1. Field of the Invention\nThe present invention relates to particular cysteine derivatives. Furthermore, the present invention relates to methods for producing particular cysteine derivatives, cosmetic agents which contain particular cysteine derivatives, and the like.\n2. Discussion of the Background\nIn human pigment cells (melanocytes) in the skin, melanin is produced utilizing L-cysteine and L-tyrosine. When a greater amount of L-tyrosine, which is a starting material for eumelanin, is utilized for the synthesis of melanin, production of eumelanin is promoted and the skin becomes dark. On the other hand, when a greater amount of L-cysteine is utilized for the synthesis of melanin, production of eumelanin is suppressed, and the skin becomes closer to yellow. Therefore, it is considered that production of eumelanin is suppressed by supplying L-cysteine during melanin synthesis.\nHeretofore, many attempts have been made to utilize L-cysteine as cosmetics such as whitening agents utilizing L-cysteine and the like. However, L-cysteine is easily oxidized, and has problems such as poor stability and bad odor for formulating as a cosmetic agent or skin external preparation.\nTo solve such problem, the development of a cysteine derivative with improved stability has been considered. JP-B-48-15938 discloses that L-2-methylthiazolidine-2,4-dicarboxylic acid is extremely stable as compared to conventional cysteine derivatives.\nIn addition, JP-A-2009-227660 discloses that a cysteine derivative obtained by esterification of L-2-methylthiazolidine-2,4-dicarboxylic acid or a salt thereof is useful as a whitening agent etc. since it has an eumelanin production suppressive effect and is stable. Furthermore, JP-A-2010-1239 discloses that 2-methylthiazolidine-2,4-dicarboxylic acid or a derivative thereof has a whitening action.\nOn the other hand, J. Biological Chem., (1937)121 539-48 discloses N-acetyl-2-methylthiazolidine-2,4-dicarboxylic acid and N-acetyl-2-methylthiazolidine-2,4-dicarboxylic acid anhydride; however, it does not describe their properties such as stability etc. and physiological activities thereof, nor does it suggest or report use for cosmetic agents or skin external preparations or whitening effects thereof."}
{"text": "Snowplows for frontal mountings on road vehicles generally are raised by the movement of a hydraulically operated piston encased in a cylinder, the moldboard with a blade on a large highway snowplow may weigh as much as 1500 to 3000 pounds.\nIt has been proposed to equip snowplow moldboards and blades with force limiters which permit the blade and moldboard to move backward and upward for cleaning obstacles (U.S. Pat. No. 3,587,182), with blade-directing and tripping devices for snowplows (e.g., U.S. Pat. No. 3,893,578), and various runners, skids shoes or caster wheels for sensing a surface by contact and, in some cases, \"planing\" a snow covering on a road.\nAlso, in the road grader art, it has been proposed to regulate the depth of cut or scraping by use of a wave energy distance sensor, of the ground, the sensor being mounted on the equipment.\nThe weight of the snowplow moldboard equipped with a blade, together with the speeds at which the vehicle-mounted snowplows frequently operate, particularly large highway plows, often makes for much blade wear and necessitates frequent blade replacement. Eight to ten hours of rugged highway service has been known to wear out a snowplow blade.\nThe present invention provides a way for reducing such blade wear and extending the operating life of the blade. In essence, it is based on regulating the fraction of moldboard and blade weight that is supported by the surface being scraped e.g., highway pavement, at a value which is effective for scraping such as snow and/or ice removal from said surface but lower then the full weight of the moldboard and blade. Such value can be, say, about 10-90% of such full weight and move frequently about 20-40% with concomitant lessening of blade abrasion."}
{"text": "1. Field of the Invention\nThe present invention relates to a front-axle bracket for motor vehicles.\n2. Description of the Related Art\nThe present invention relates to a front-axle bracket for motor vehicles.\nSuch brackets, which are also referred to as saddles, frames, or auxiliary frames, are attached to vehicle structures, for example to the car body and/or to longitudinal supports of a vehicle, as pre-assembled units or modules, after they have been provided with assemblies or auxiliary assemblies.\nA front-axle bracket is disclosed, for example, in WO 2007/031060 A1, which goes back to the same applicant."}
{"text": "1. Field of the Invention\nThe present invention generally relates to an electronic control unit of a kind required to have a waterproof structure suitable for use in an automobile and, more particularly, to the electronic control unit equipped with at least one electrical input/output connector having a plurality of terminal members adapted to be connected directly with a printed circuit board.\n2. Description of the Prior Art\nThe electrical plug-in connector assembly providing for the input and output of one or both of an electrical power and electrical signals to the separate device has long been well known, which is generally made up of a plug-in connector and a socket connector for receiving the plug-in connector. While the plug-in connector assembly is available in various types, the present invention makes use of the type comprising the plug-in connector having a plurality of terminal sheaths and the socket connector having a corresponding number of connector terminals which are connected at one end to a printed circuit board by means of soldering and which are, when the plug-in connector is plugged into the socket connector, engaged or plugged into the associated terminal sheaths to establish electric circuits between the connector terminals and the terminal sheaths.\nIn this particular type of the plug-in connector assembly, it is well understood that whenever the plug-in connector is plugged in or removed from the socket connector, a compressive or tensile load act on the printed circuit board, respectively. Where the socket connector is fixed in position with the connector terminals electrically connected directly with the printed circuit board through electrical joints formed by, for example, soldering, repeated application of the alternating compressive and tensile loads to the printed circuit board will eventually result in degradation of the reliability of the printed circuit board. Accordingly, it is a general practice to mount the plug-in connector and the printed circuit board separately on a housing of the electronic control unit in order to secure a highly reliable electronic control unit.\nIn the structure in which the socket connector and the printed circuit board are fixed separately in the housing of the electronic control unit, the tensile load which would be brought about when the plug-in connector is removed from the socket connector will not be transmitted to the printed circuit board. However, the difference between the coefficient of thermal expansion of a material for the connector terminals in the socket connector and that of a material for the printed circuit board is apt to allow thermal stresses to develop at the solder joints which eventually leads to material fatigue at the solder joints. Once this occurs, breakage is likely to result in at one or some of the solder joints.\nFIG. 4 illustrates one of the conventional methods of accomplishing an electrical connection between one of the connector terminals and the printed circuit board. As shown therein, a control unit housing 100 is made of a synthetic resin and has a connector casing 101 and a circuit casing 103 formed integrally therewith with a partition wall intervening therebetween. The connector casing 101 has a socket defined therein so as to open upwardly for receiving a plug-in connector not shown. The illustrated socket connector also includes a plurality of straight connector terminals, only one of which is shown by 102, that are fixed in position by interference fit or insert-molding so as to extend through the partition wall of the control unit housing 100. The straight connector terminals 102 thus have one end positioned inside the socket in the connector casing 101 and the opposite end positioned inside the circuit casing 103. A printed circuit board 105 is housed within and secured to the circuit casing 103 by means of a plurality of set screws, only one of which is shown by 104, that are screwed into an inwardly protruding peripheral portion 106 integral with the circuit casing 103. Those ends of the straight connector terminals 102 within the circuit casing 103 are pierced through corresponding holes in the printed circuit board 105 and are then soldered at 102a to circuit elements on the printed circuit board 105.\nAccording to the connection shown in FIG. 4, the connector terminals 102 are simple in shape and can easily and accurately soldered to the printed circuit board 105. However, considering that the control unit housing 100 including the circuit casing 103 is made of the synthetic resin whereas the connector terminals 102 are made of, for example, brass having a good conductivity, each time the ambient temperature inside a space h delimited between the wall of the circuit casing 103, a part of which concurrently serves as the partition wall between the connector casing 101 and the circuit casing 103, and the printed circuit board 105 undergoes change, the solder joints 102a tends to be pulled laterally due to the peripheral portion 106 and, hence, the circuit casing 103 having a larger coefficient of thermal expansion, which would eventually results in failure of some or all of the solder joints 102a to conduct an electrical current therethrough.\nIn order to avoid the problem inherent in the structure shown in and described with reference to FIG. 4, an attempt has largely been made to provide each connector terminal of the socket connector with a deformable portion. By way of example, as shown in FIG. 51 the use is made of a connector subassembly 200 comprising a connector casing 201, separate from and independent of the circuit casing 103, and a plurality of connector terminals 202, that are fixed in position by interference fit or insert-molding so as to extend through a wall of the connector casing 201. The connector terminals 202 are so bent into a generally L-shaped configuration as to define the deformable portion so that in the assembled condition, neither the compressive load nor the tensile load will be transmitted to the solder joints 202a.\nThe connector subassembly 200 of the structure described above is fixedly secured to the circuit casing 103 with the connector terminals 202 extending into the circuit casing 103 through one or more perforations defined in a wall of the circuit casing 103 and then soldered at 202a to the printed circuit board 105 in a manner similar to that described in connection with the structure of FIG. 4. As is the case with the structure shown in FIG. 4, the printed circuit board 105 is supported within the circuit casing 103 by means of a plurality of set screws 104 screwed into the inwardly protruding peripheral portion 106 integral with the circuit casing 103.\nIn the structure shown in and described with reference to FIG. 5, change in dimension due to the difference in coefficient of thermal expansion between the connector terminals 202 and the control unit housing 100, which change tends to be brought about by change in ambient temperature inside the circuit housing 103, can be effectively absorbed by the resiliency of the deformable portion of each connector terminal 202. Consequently, no external force can be transmitted to the solder joint 202a to thereby avoid the possibility of electrical disconnection between the connector terminals 202 and the corresponding circuit elements on the printed circuit board 105.\nTo complete the control unit housing 100, however, the structure shown in FIG. 5 requires the use of the separate casings 103 and 201 which tends to increase the cost of manufacture. In addition, where the waterproof design is required, the use of the separate casings 103 and 201 requires a seal to be employed between these casings. Accordingly, not only does the cost of manufacture increase additionally, but also the use of the waterproof seal may constitute a cause of reduction in reliability. Where the fluid tightness is of an paramount importance, the electronic control unit should have the connector casing as a unitary part of the control unit housing 100 together with the circuit casing 103.\nHowever, another problem will arise if the connector casing 201 is formed integrally with the control unit housing 100 together with the circuit casing 103, which will now be discussed with reference to FIG. 6.\nThe control unit housing 100 shown in FIG. 6 has the partition wall dividing it into the connector casing 101 and the circuit casing 103, both of which is a unitary part of the control unit housing 100. When it comes to press-fitting of the L-shaped connector terminals 202 across the partition wall, the socket in the connector casing 101 has such a limited capacity that each connector terminal 202 cannot be passed through the partition wall with its opposite ends positioned on respective sides of the partition wall. Accordingly, it would be a reasonable way to insert each connector terminal 202 from the interior of the circuit casing 103 instead of from the socket in the connector casing 201, in which case a wall portion of the circuit casing 102 confronting the partition wall must be set back a distance, indicated by L2, to provide an extra space to accommodate an end portion of each connector terminal 202 that is passed through the partition wall. The set-back distance L2 is as a matter of design equal to or greater than the length L1 of that end portion of the respective connector terminal 202, rendering the control unit housing 100 as a whole to be bulky in size.\nIn addition, alignment of the tip of that end portion of each connector terminal 202 with a corresponding perforation in the partition wall is indeed difficult to achieve and even the slight misalignment would require readjustment in position of the connector terminals 202 relative to the respective perforations in the partition wall. This complexity indeed increases the cost of manufacture.\nAlthough the connector terminals 202 may be insert-molded together with the control unit housing 100 during the molding of the latter, a complicated and expensive mold would be required in such case, adding cost to the final product.\nAlso, considering a top surface area of the circuit casing 103 on one side of the partition wall opposite to the connector casing 201 defines a support base on which a hydraulic unit (not shown) including a hydraulic pump and a plurality of electronically controlled control valves is fixedly mounted, the provision of the set-back distance L2 results in a considerable reduction in efficiency of maximum utilization of the available space. In addition, since the structure shown in FIG. 6 requires the connector terminals 202 to be press-fitted into the corresponding apertures in the partition wall from the interior of the circuit casing 103, it is impossible to reduce the size of a portion where the connector terminals 202 are fixed even though they can be formed to have the same shape as the corresponding socket terminals in the plug-in connector, and a problem associated with reduction in physical strength of the partition wall would also occur if the connector terminals 202 are spaced a small pitch between the neighboring members of them.\nU.S. Pat. No. 5,452,948, issued Sep. 26, 1995 to Cooper et al. discloses the electronically controlled hydraulic unit comprising, as reproduced in FIG. 6, the idea of bending coil leads or terminals, extending respectively from each of a plurality of electromagnetic coils, to provide resiliency so that the solder joints between circuit elements on a coil printed circuit board and the coil terminals can be substantially immune from being stressed when a compressive load is applied to the coil in a direction parallel to the coil terminals.\nMore specifically, referring to FIG. 7, Cooper et al. U.S. Patent discloses an electronically controlled hydraulic unit comprising a coil mounting housing 300 including a peripheral side wall 301 and a connector casing 302 both formed integrally therewith. Electromagnetic coils 303 each having a pair of coil terminals 306 are each seated on a respective coil mounting seat 304 having a snap-in, hold-down projection 305. The coil terminals 306 of each of the electromagnetic coils 303 are soldered at a free end thereof to circuit elements on a coil printed circuit board 307, but have a generally intermediate portion thereof bent to provide resiliency so that the solder joint will not be stressed when a compressive load is applied to the respective coil 303 in the direction parallel to the coil terminals 306 during mounting of the respective coil 303 onto the coil mounting housing 300.\nThe coil mounting housing 300 has a plurality of staking pins 308 protruding upright from a sealing rim extending inwardly from the bottom of the peripheral side wall 301. These staking pins 308 are, after the coil printed circuit board 307 has been mounted in the coil mounting housing 300 with its peripheral portion resting on the sealing rim through a sealing member 309, bent over to hold the printed circuit board securely in place. An electronic control unit printed circuit board 311 having a number of control elements 31 mounted thereon is positioned beneath the coil mounting housing 300 and is supported by connecting legs 312, which extend downwardly from a socket connector 302 and are then soldered to the electronic control unit printed circuit board 311, and also by pins 313 having one end fixed to the coil printed circuit board 307 so as to extend downwardly therefrom and soldered to the electronic control unit printed circuit board 311. After the component parts have been assembled, a lower cover 314 is secured from below to the coil mounting housing 300 to thereby maintain sealing integrity between the coil printed circuit board 307 and the coil mounting housing 300.\nIn this known electronically controlled hydraulic unit, the coil printed circuit board 307 is electrically connected with the electromagnetic coils 303 through the coil terminals 306 and is securely connected to the coil mounting housing 300 through the sealing member 309. On the other hand, the electronic control unit printed circuit board 311 are electrically and mechanically connected with the terminal legs 312 and the pins 313, respectively. Since the lower cover 314 when secured from below to the coil mounting housing 30 encloses the electronic control unit printed circuit board 311 in a substantially hermetic fashion, the various component parts can be fluid-tightly sealed.\nAccording to Cooper et al. U.S. Patent, the coil terminals so bent as to provide resiliency appear effective to avoid the possibility of the solder joints from being stressed. However, the bending of the coil terminals for this purpose is required either before or after the coil terminals have been pierced under interference fit through the coil mounting seats. In either case, it has been found that some of the coil terminals are apt to deviate from the design dimension when they are bent and, therefore, alignment of the coil terminals with corresponding perforations in the coil printed circuit board appears to be difficult to achieve without complicated and time-consuming procedures.\nIn general, soldering is the most inexpensive and effective means for connecting connector terminals in an electrical connector with a printed circuit on a printed circuit board. However, the resultant solder joints have been found sensitive to external stresses such as, for example, thermal stresses generated in the circuit board or the like which would be brought about by change in ambient temperature."}
{"text": "Life safety systems include notification devices positioned on walls and ceilings, as required by law. These notification devices can include a horn for an audible alert and/or a strobe for a visual alert. Notification devices draw energy to illuminate a strobe in a life safety application, such as a fire alarm or other signaling device. A notification device can be mounted on a surface, e.g., a wall or a ceiling, and aimed at the floor. The notification device has a reflector with internal reflecting surfaces to reflect light toward the specified space. The reflector can also have some external secondary reflector elements mounted on the reflector body.\nNotification devices for the hearing impaired are governed by standards that require a polar light distribution off-axis at any given rated candela. In order to achieve the standard, the notification device must produce a minimum output at specified angles off-axis. For example, for a notification device to be rated at 185 candela, UL 1971 requires that the strobe has specific light intensities at viewing angles ranging from zero to 90 degrees off axis. UL 1971 requires a polar light distribution pattern to enhance the likelihood of alerting hearing impaired individuals throughout an area. The light intensity is measured in a horizontal and vertical direction and at viewing angles ranging from zero to +90 degrees and zero to −90 degrees. In one example, the UL polar distribution requires a percentage of the candela output at specific angles off-axis, e.g., five degrees off-axis requires an output of 90% of the rated on-axis value. So a 15 candela strobe is required to output 13.5 candela at five degrees off-axis.\nIt is desirable to use less energy, but the light output of the strobe must meet the requisite standards to achieve underwriting. As a result, the amount of energy is limited by the minimum required light output. The amount of current drawn by the notification device can be affected by various aspects of its configuration, including reflector design, electrical efficiency, lamp efficiency, efficiency of a metalized coating used for high reflectance, and the efficiency of the lens. It is desirable to adjust the properties of the notification device so that less current is needed to power the notification device while maintaining a required amount of light output."}
{"text": "The present invention relates to a phase shifter used in a demodulator of a digital-signal-modulated carrier wave transmission system.\nIn a demodulator for demodulating a digital-signal-modulated carrier wave, a modulated signal input to the demodulator is supplied to two phase detectors, and output signals having predetermined phase angles shifted through, e.g., 90.degree. from each other are supplied from the two phase detectors. Therefore, a reference carrier signal from a VCO or the like is directly supplied to one phase detector and to the other phase detector through a 90.degree.-phase shifter. A control signal is supplied to the phase shifter so that the output from the phase shifter always accurately maintains a phase difference of 90.degree. with respect to the output from the VCO.\nA conventional phase shifter used in a digital demodulator is disclosed in U.S. Pat. No. 4,549,142. This phase shifter is arranged using an LC circuit and therefore has frequency dependency. For this reason, when this phase shifter is mounted on a printed circuit board or the like, a phase shift is easily adversely affected by a parts arrangement. Therefore, this phase shifter is not suitable for an LSI arrangement.\nJapanese Patent Laid-Open No. 61-157019 discloses various phase shifters using not an inductance but a capacitor, a resistor, and a transistor, thereby eliminating the above drawback. However, a phase shift amount cannot vary in these phase shifters."}
{"text": "1. Field of the Invention\nThis invention is directed to an electronic ballast system for fluorescent or gas discharge tubes. In particular, this invention relates to automatic gain control ballast systems for fluorescent tubes. More in particular, this invention relates to a ballast system which provides for frequency control mechanisms utilizing inductance characteristics of an inverter transformer to allow for a substantial frequency stabilization. Still further, this reference directs itself to a transistor switching mechanism which is feedback coupled to an induction circuit for establishing a regulated current for a gas discharge tube independent of the gain of a transistor used in a particular ballast system unit. Further, this invention directs itself to an electronic ballast system where there is provided a transistor switching network including a regulation control circuit for maintaining the power output of a gas discharge tube at a predetermined and substantially constant value. Still further, this invention relates to an electronic ballast system having a toroidal transformer to provide a predetermined variable inductance for regulating a power output to a gas discharge or fluorescent tube. Additionally, this invention pertains to an electronic ballast system which is transistorized and where the current gain of various transistors range over a wide value range from one system unit to another and provides for electrical circuitry which will maintain the gas discharge tube light output fluctuation to a minimum value.\n2. Prior Art\nElectronic ballast systems for gas discharge or fluorescent tubes are known in the art. However, in some prior art electronic ballast systems, no provision is made for frequency stabilization of the circuit. Thus, in such prior art electronic ballast system circuits, when a gas discharge tube is removed from the circuit, there is a deleterious flickering of the remaining gas discharge tubes, or in some cases, a complete breakdown of the visible light from the remaining gas discharge tubes.\nIn other prior art electronic ballast systems, the light output of the gas discharge tubes are highly dependent upon the gain of the transistors used in the circuits. In such prior art where the transistor gains between one unit and another vary over a great range, the light output from the gas discharge tube fluctuates by large amounts. In such prior art systems, additional circuitry must be added to such prior art electronic ballast systems to maintain the light output fluctuation as constant as possible between different units.\nIn the alternative, in such prior art electronic ballast systems, in order to maintain the light output fluctuation of the gas discharge tubes to a minimum value, the transistors used in the circuit must be matched within very restrictive ranges to obtain a substantially constant transistor gain from one unit to the next. Such procedures increase the cost of the electronic ballast systems, as well as providing for additional electronic elements, which increase both labor and capital investment in such electronic ballast systems.\nBy increasing the electrical components used in prior art electronic ballast systems to maintain the gas discharge tubes in operation, there is a responsive increase in the manufacturing costs of such prior art electronic ballast systems and a reduction in the reliability of these ballast systems."}
{"text": "The security of computing resources and associated data is of high importance in many contexts. As an example, organizations often utilize networks of computing devices to provide a robust set of services to their users. Networks often span multiple geographic boundaries and often connect with other networks. An organization, for example, may support its operations using both internal networks of computing resources and computing resources managed by others. Computers of the organization, for instance, may communicate with computers of other organizations to access and/or provide data while using services of another organization. In many instances, organizations configure and operate remote networks using hardware managed by other organizations, thereby reducing infrastructure costs and achieving other advantages. With such configurations of computing resources, ensuring that access to the resources and the data they hold is secure can be challenging, especially as the size and complexity of such configurations grow."}
{"text": "My U.S. Pat. No. 4,367,850 describes a cartridge including a length of magnetic tape spliced into an endless loop. In one embodiment that cartridge can be made to conform to the standard of the NAB so that it can be used in the types of record/playback machines presently in use in the broadcast industry. That cartridge embodiment causes much less edge wear on the tape and much less change in tension in the tape at a head of a machine in which the cartridge is received than the other types of cartridges presently in use in the broadcast industry.\nThe cartridge described in U.S. Pat. No. 4,367,850 employs a method for maintaining a uniform high tension in the tape at a head of a machine in which the cartridge is engaged so that the tape can be pressed against the head by tension alone without the use of pressure pads. That method comprises providing a fixed, generally cylindrical hub having a central opening and a slot extending axially across the full width of the hub and communicating with its central opening, and an endless length of magnetic tape; wrapping a major portion of the tape about the hub to form a coil while allowing a minor portion of the tape to extend from the innermost wrap of the coil through the slot into the central opening of the hub, and around the side surface of the coil to the outermost wrap of the coil; pulling the tape from the slot and across a head on a record/playback machine; and applying a light force to tension the minor portion of the tape as it moves onto the outer wrap of the coil which will produce the high, generally uniform tension in the minor portion of the tape being pulled from the coil and across the head.\nThe embodiment of the cartridge adapted for use in the broadcast industry which employs that method to tension tape across a head comprises a housing adapted to be received in a record/playback machine and having access openings adapted to receive record/playback heads and a tape drive mechanism in the machine. The hub is fixed on the housing at a position spaced from the access openings. Means on the housing define a tape path for, and produce tension in, the minor portion of the tape. Those means comprise means for guiding the minor portion of the tape past the access openings in a predetermined position for engagement by the heads and the drive mechanism of the playback machine, and a guide pin guiding the minor portion of the tape between the access openings and the outermost wrap of the coil. The guide pin is mounted for movement between a first position providing a first path length between the access openings and the outer wrap of the coil and a second position providing a second path length longer than the first path length between the access openings and the outer wrap of the coil (which means preferably is an arm having a first end supporting the guide pin, and a second end pivotably mounted on the housing to afford movement of the pin along an arcuate path adjacent the periphery of the coil between its first and second positions), and means are provided for biasing the guide pin toward its second position.\nThe pin can move to positions between its first and second positions under the influence of the biasing means to accommodate changes in length of the minor portion of the tape which decreases or increases respectively when the length of the major portion increases or decreases. The major portion of the tape cyclically undergoes its largest change in length by slowly decreasing in length as the splice moves from the coils outer wrap toward its inner wrap and by then suddenly increasing in length as the splice leaves the coil. The means for biasing the guide pin is adapted to apply a small force at the guide pin (e.g., generally in the range of 2 to 14 grams) to lightly tension the tape extending around the pin and moving onto the outermost wrap of the coil, somewhat in the manner of a rope or Proney brake, which light tension produces a significantly higher tension (e.g., generally in the range of 50 to 115 grams or 2 to 4 ounces) in tape leaving the coil. That higher tension is surprisingly uniform despite small changes in the force applied by the guide pin as the length of the minor portion of the tape changes due to the position of the splice along the tape.\nThe tape is guided so that the quite uniform higher tension thus produced in the minor portion of the tape between the inner wrap of the coil and the drive mechanism of a machine in which the cartridge is engaged presses the tape against the record and playback heads of the machine with sufficient pressure that pressure pads are not required, and the tape is thus not subjected to the abrasion and erratic forces caused by sliding contact between the tape and such pressure pads.\nAlso, my U.S. Pat. No. 4,394,989 describes a tape cartridge of the type described above in which the length of the minor portion of the tape can be easily adjusted by persons in the field to accommodate changes in the cartridge such as tape wear to ensure that a desired range of guide pin movement will be retained.\nIn that improved cartridge, the hub comprises a flexible arcuate cantilevered portion having a first end partially defining the slot, a second end that is fixed on the housing, and which is separated from the housing between its ends. Means in the form of a cam rotatably mounted on the housing is provided for changing the position of the cantilevered hub portion radially of the rest of the hub to adjust the effective diameter of the hub and thereby the length of the minor portion of the tape to cause the guide pin to move within a desired range of movement between its first and second positions. The cam has a peripheral surface contacting the inner side of the cantilevered hub portion, and can be manually rotated with a tool such as a screwdriver as the tape is being propelled by a machine so that the effects of cam adjustments on the tape can be observed as they are made. Such adjustability ensures that changes in overall tape length which can occur due to causes such as tape wear can be compensated for. While such adjustments do not need to be made often during the life of a cartridge, they do require attention by technicians, such as each time the tape is re-recorded.\nThus the cartridge described in my U.S. Pat. No. 4,394,989 is not particularly useful in cartridges used in background or foreground music machines where users typically are not trained to or interested in making the occasional cam adjustments needed to maintain the cartridge in its best operating condition.\nAnother problem that occurs with the cartridge described in U.S. Pat. No. 4,394,989 is that when it is moved from a room at normal temperatures into the cold (e.g., from an area at about 68.degree. F. to an area at below 10.degree. F.) its coil and fixed hub will shrink. When the cartridge is then again brought into a warm room and playing of the cartridge is attempted, its tape coil will seize around the hub which is more exposed and returns to room temperature and thus expands to its original size more quickly than the tape coil. Radio technicians can correct for this problem by adjusting the cam to release the tension in the coil. Again, however, persons involved with background music or foreground machines typically lack the training and desire to perform similar adjustments."}
{"text": "1. Field of the Invention\nThe present invention relates to the field of power supply, and more particularly, to a power supply apparatus for an LED lamp, which utilizes the principle of persistence of vision of human eyes to modulate the waveforms and frequencies of voltages for driving LEDs and can still maintain the normal operation.\n2. Description of the Related Art\nAs compared with conventional light sources, light emitting diodes (LEDs) have the advantages of operating at a low voltage, low power consumption and a long service life, etc. Therefore, in modern life, LEDs have been widely used in various related fields, where a light source is required, for example, applied to display backlight modules, indicating light sources, general lighting equipment, and the like. LEDs have the especially significant advantage of energy saving, so their practicability is particularly important in the most general purpose of lighting.\nThe principle of a traditional power supply for driving an LED lamp is mainly that an inputted high voltage alternating current is rectified, filtered and then converted by a transformer into a low voltage alternating current then being processed by a rectifier and an electrolytic capacitor filter circuit into a stable voltage DC output for driving corresponding LED lamps. The prior art uses large capacitance characteristics of electrolytic capacitors to reduce the occurrence of ripples and thus to achieve a stable voltage DC output. However, most external environmental factors will deteriorate the performance of electrolytic capacitors such that electrolytic capacitors typically have a shorter life than other types of capacitors. Among environmental factors, temperature has the most tremendous impact on the life of electrolytic capacitors, and for example, the intense heat caused by ripple currents is one of the reasons. The intense heat will accelerate the performance deterioration of electrolytic capacitors to reduce the life and static capacity to only a fraction of the original, and thus the capacitors are tantamount to having no filter capacitor function. Therefore, even if an LED lamp has a longer life, there is no filter capacitor to effectively match the longer life for the normal operation of the LED lamp.\nIn view of the above-described circumstances, the invertors conducted elaborate research with accumulated years of experience in this field, so as to develop a power supply apparatus for an LED lamp, which mainly utilizes the principle of persistence of vision of human eyes to modulate the waveforms and frequencies of voltages for driving LEDs and can still maintain the normal operation. Also, with no need to convert an AC voltage into a stable DC output as in the prior art, the present invention uses a longer life non-electrolytic capacitor for wave filtering to control the waveforms and frequencies of any voltages outputted to an LED lamp to be greater than or equal to 60 Hertz and control peak-to-peak values of ripple signals to be greater than or equal to 1 volt. Furthermore, feedback control facilitates to maintain the output within the operating range so as to effectively drive the LED lamp, thus improving the overall service life of the LED lamp."}
{"text": "The concept of wire connectors for connecting the junction of two or more wires together by twisting a cap on the wires is old in the art. Wire connectors are well known in the art and generally comprise an outer housing with a tapered threaded interior to permit a user to insert wires into the tapered opening. To use a wire connector, the user inserts the twisted ends of electrical wires into a cavity on the inside of the wire connector. The user then holds the wire in one hand and with the other hand twists the wire connector. The twisting action pulls the junction ends of the wires into a low resistance electrical contact.\nIf the connector is located in a wet location it is necessary to place a waterproof sealant around the connector. In order to prevent water or moisture from entering the connector and forming an oxidation layer over the ends of the wire the user inserts the wire connector and the wire into some type of a waterproof potting compound. The compound may be either a non hardening or a hardening compound. In either case the compound creates a waterproof capsule over the wire connector and the junction ends of the electrical wires.\nThe prior art process is time consuming because it involves two separate steps as well as the nuisance of having separate potting compounds and containers to hold the potting compound. A second generation improved twist-on wire connector wherein the wires can be encapsulated and sealed in a twist on wire connector to prevent water or moisture from entering the connector is shown in my U.S. Pat. Nos. 5,113,037; 5,023,402 and 5,151,239. The second generation twist-on wire connectors permits the user in one continuous action to simultaneously form the junction ends of wire leads into a low resistance electrical connection that is surrounded by a waterproof sealant to form a waterproof covering around the junction ends of the wire leads.\nUnder certain dynamic conditions, such as vibration and shock, or large temperature changes the wires in the twist-on wire connector can become loosened and thus lower the integrity of the connection between the wires in the twist on wire connectors by either increasing the electrical resistance or decreasing the contact area or both. I call the present invention a third generation electrical twist-on wire connector wherein the integrity of the low resistance electrical connection of the twist-on wire connector is actually enhanced by placing a small amount of self-adhering lubricant in the twist-on wire connector. Generally, to enhance the electrical conductive between connector and wire one needs only a small amount of self-adhering lubricant to provide an enhanced low resistance electrical connection. In the present invention, it has been found that when a small amount of a self-adhering lubricant has been incorporated into the wire connector it results in an enhanced low resistance electrical connection.\nIn order to ensure that twist-on wire connectors remain in a low-resistance electrical connection when subject to environmental conditions of shock and vibration one prior art method is to encapsulate the twist-on wire connectors in a solid resin. In other methods tools may be used to apply extra torque to the wires to cause further engagement of the threads of the wire connector with the electrical wires. In the present method one need not resort to encapsulation with a solidified resin and one need not resort to using tools to enhance the torque since it has been found that even though a lubricant is used the suspected detrimental effects of use of a lubricant film between the wire and the wire connector is overcome by the greater physical engagement between the wires which can be obtained by use of the lubricant. That is, the lubricant allows one to reduce the torque that may be required to make a secure connection in a twist-on wire connector. As a result one can hand tighten a twist-on wire connector containing a small amount of self-adhering lubricant and provide for greater electrical contact between the wires which results in retention of the low-resistance electrical connection even in the presence of forces such as vibration and shock."}
{"text": "1. Field of the Invention\nThe present invention relates generally to the production of electric power by use of geothermal fluids and more particularly to processes for reducing the silica content of flashed, silica-rich geothermal brine by inducing crystallization of the silica from the brine onto seed crystals.\n2. Discussion of the Prior Art\nLarge subterranean aquifers of naturally produced (geothermal) steam or hot aqueous liquids, specifically water or brine, are found throughout the world. These aquifers, which often have vast amounts of energy potential, are most commonly found where the earth's near-surface thermal gradient is abnormally high, as evidenced by unusually great volcanic, fumarole or geyser activity. Thus, as an example, geothermal aquifers are fairly common along the rim of the Pacific Ocean, long known for its volcanic activity.\nGeothermal steam or water has, in some regions of the world, been used for centuries for therapeutic treatment of physical infirmities and diseases. In other regions, such geothermal fluids have long been used to heat dwellings and in industrial processes. Although efforts to further develop geothermal resources for these site-restrictive uses continue, considerable recent research and development has, instead, been directed to exploitation of geothermal resources for production of electrical power which can be conducted, often over existing power grids, for long distances from the geothermal sources. In particular, recent steep increases in the cost of petroleum products used for conventional production of electric power, as well as actual or threatened petroleum fuel shortages or embargos have intensified the interest in use of geothermal fluids as an alternative and generally self-renewing source of power plant \"fuel\".\nGeneral processes by which geothermal fluids can be used to generate electric power are known and have been known for some time. As an example, geothermal steam, after removal of particulate matter and polluting gases, such as hydrogen sulfide and ammonia, can be used in the manner of boiler-generated steam to operate steam turbine generators.\nNaturally pressurized geothermal brine or water having a temperature of over about 400.degree. F. can be flashed to a reduced pressure to convert some of the brine or water to steam. The steam produced in this manner can then be used to drive steam turbine generators. The flashed geothermal liquid and the steam condensate obtained from power generation are typically reinjected into the ground to replenish the aquifer and prevent ground subsidence. Cooler geothermal brine or water can often be used to advantage in binary systems in which a low-boiling point, secondary liquid is vaporized by the hot geothermal liquid, the vapor produced being used to operate gas turbine generators. The cooled brine is typically reinjected into the ground.\nAs might be expected, use of geothermal steam is preferred over use of geothermal water or brine for generating electric power because the steam can be used more directly, easily and cheaply. Consequently, where readily and abundantly available, geothermal steam has been used for a number of years to generate commercially important amounts of electric power at favorable costs. For example, by the late 1970's, geothermal steam at The Geysers in Northern California was generating about two percent of all the electricity used in California.\nWhile energy production facilities at important geothermal steam sources, such as at The Geysers, are still being expanded, when not already at capacity, the known number of important geothermal steam aquifers is small compared to that of geothermal brine or water. Current estimates are, in fact, that good geothermal brine or water sources are about five times more prevalent than are good sources of geothermal steam. The potential for generating electric power is, therefore, much greater for geothermal brine and water than it is for geothermal steam. As a result, considerable current geothermal research is understandably directed towards the development of economical geothermal brine and water electric generating plants, much of this effort being expended towards use of vast geothermal brine resources in the Imperial Valley of southern California.\nAlthough, as above mentioned, general processes are known for using geothermal brine or water for production of electric power, serious problems, especially with the use of highly saline geothermal brine, have often been encountered in practice. These problems have frequently been so great as to prevent the production of electric power at competitive rates and, as a consequence, have greatly impeded the progress of flashed geothermal brine power plant development in many areas.\nThese severe problems are caused primarily by the typically complex composition of geothermal brines. At natural aquifer temperatures in excess of about 400.degree. F. and pressures in the typical range of 400 to 500 psig, the brine leaches large amounts of salts, minerals and elements from the aquifer formation, the brine presumably being in chemical equilibrium with the formation. Thus, although brine composition may vary from aquifer to aquifer, wellhead brine typically contains very high levels of dissolved silica, as well as substantial levels of dissolved heavy metals such as lead, copper, zinc, iron and cadmium. In addition, many other impurities, particulate matter and dissolved gases are present in most geothermal brines.\nAs natural brine pressure and temperature are substantially reduced in power plant steam conversion (flashing) stages, chemical equilibrium of the brine is disturbed and saturation levels of impurities in the brine are typically exceeded. This causes the impurities and silica to precipitate from the brine, as a tough scale, onto surrounding equipment walls and in reinjection wells, often at a rate of several inches in thickness per month. Assuming, as is common, that the brine is saturated with silica at the wellhead, in high temperature portions of the brine handling system, for example, in the high pressure brine flashing vessels, heavy metal sulfide and silicate scaling typically predominates. In lower temperature portions of the system, for example, in atmospheric flashing vessels, amorphous silica and hydrated ferric oxide scaling has been found to predominate. Scale, so formed, typically comprises iron-rich silicates, and is usually very difficult, costly and time consuming to remove from equipment. Because of the fast growing scale rates, extensive facility down time for descaling operations may, unless scale reducing processes are used, be required. Associated injection wells may also require frequent and extensive rework and new injection wells may, from time to time, have to be drilled at great cost.\nTherefore, considerable effort has been, and is being, directed towards developing effective processes for eliminating, or at least very substantially reducing, silica scaling in flashed geothermal brine handling systems. One such scale reduction disclosed in U.S. Pat. No. 4,370,858 to Awerbuch, et al, involves the induced precipitation of scale-forming materials, notably silica, from the brine in the flashing stage by contacting the flashed brine with silica or silica-rich seed crystals. When the silica saturation level in the brine is exceeded by the brine being flashed to a reduced pressure, supersaturation amounts of the silica leaving solution in the brine deposit onto the seed crystals. Not only do the vast number of micron-sized seed crystals introduced into the flashing stage provide a very much larger surface area than the exposed surfaces of the flashing vessels but the silica from the brine tends to preferentially deposit onto the seed crystals. Substantially all of the silica supersaturation, therefore, precipitates onto the seed crystals instead of precipitating as scale onto vessel and equipment walls and in injection wells.\nPreferably, the seed crystals are introduced into the high pressure flashing vessel, which may then be referred to as a high pressure flash crystallizer, wherein the brine first becomes supersaturated in scale-forming materials. The crystallization process, while starting in the high pressure flash crystallizer, continues in successive, lower pressure flashing vessels in which the brine typically again becomes supersaturated with silica. In a downstream reactor-clarifier, the silicious precipitate is separated from the brine as a slurry which contains about 30 percent by weight of silica. According to known processes, a portion of this silicious slurry from the reactor-clarifier stage is recirculated back upstream into the high pressure flash crystallizer, whereby the silica material in the slurry acts as seed material.\nAfter subsequent filtering to remove fine silicious particles not removed in the reactor-clarifier stage, the \"clarified\" brine is commonly reinjected into the ground in an injection stage.\nAs above-mentioned, geothermal brines used for electric power generation are, at wellhead temperature and pressure, frequently saturated with silica. As a consequence, substantial amounts of silica must be precipitated from the brine onto the seed material in the flash crystallization stage in order to prevent silica scaling in downstream brine handling equipment. Such removal of silica from the geothermal brine requires, particularly for high brine flow rates associated with production of reasonably large amounts of power, effective and rapid silica precipitation so that brine residence time in the flash crystallizer vessels, as determined by vessel capacity, can be maintained within acceptable and practical limits.\nSuch known silica seeding processes which use the silicious slurry from reactor-clarifier stage as seed material have, however, substantial disadvantages. A major disadvantage is that because of the large volume of seed material required for rapid, effective silica precipitation and because of the lower temperature of the seed slurry compared with that of the main flow of brine into which the slurry is introduced, steam production in the flash crystallizer is significantly reduced by quenching action of the slurry. The amount of steam \"lost\" as a result of such quenching may be about 5 percent. For an exemplary 10 megawatt power plant a 5 percent steam loss is equivalent to about 0.5 megawatts power loss, having a current value of about $33 an hour or about a quarter of a million dollars for a typical, 10 month operational year. Another disadvantage is that substantial costs are associated with the purchase, installation, operation and maintenance of the pipes, fittings and relatively large pumps required for pumping the seed slurry from the reactor-clarifier stage upstream to the high pressure flash crystallizer.\nIt is therefore, an object of the present invention to provide a process for removing silica from flashed geothermal brine by in-situ seed material formation in the flash crystallization stage of a geothermal brine handling system.\nAnother object of the present invention is to provide a process for removing silica from flashed geothermal brine by in-situ seed material formation in the flash crystallization stage of a geothermal brine handling system by reacting a base with naturally present heavy metals dissolved in the flashed brine so as to form insoluble heavy metal compounds which function as seed material.\nOther objects, advantages and features of the present invention will become apparent to those skilled in the art from the following description, when taken in conjunction with the accompanying drawing."}
{"text": "This invention relates generally to fluid dispensers, and specifically to a class of fluid dispensers adapted for attachment to fluid containers such as flexible plastic bags.\nThe use of plastic bags as fluid containers has become widespread, particularly in certain applications. For example, food liquids, such as milk and wine, are commonly packaged and distributed in plastic bags. The filled bags are usually carried in a cardboard box. A rigid plastic dispenser is attached to the bag in a manner to extend through the box in order to dispense liquid from the bag. The dispenser normally includes a valve which is operated between open and closed positions by the ultimate user of the product, in order to periodically remove liquid from the bag. As liquid is removed from the flexible plastic bag, the bag collapses around the remaining liquid. Since it is unnecessary that air be introduced into the container in order to remove liquid, air does not come into contact with the remaining food material in the bag.\nDuring shipping and storage, the bag remains sealed. One type of commercial dispenser includes a cutting element for making an opening in the bag wall. The first time the dispenser is operated to withdraw fluid from the bag, the cutting element cuts a hole in the bag wall. Prior to this time, the bag has remained totally sealed. A manually operated valve provided as part of the dispenser thereafter controls flow of the fluid out of the bag through the bag opening. Such an arrangement is described in U.S. Pat. No. 4,355,737 Pongrass and Rutter (1982). When this type of dispenser is used, the bag is initially filled through an opening that is heat sealed shut after filling.\nAnother type of commercial dispenser is attached to a bag that includes a hole pre-punched in its wall. Before attaching the dispenser to the bag, it is filled through the hole with liquid, particulate material, or other fluid desired to be stored and shipped in it. A rigid plastic spout is usually first attached to the bag around the hole. A dispenser is then attached to the spout and seals the bag opening by a diaphragm extending across the fluid passage within the dispenser. This diaphragm is then cut when the user first opens the dispenser valve to remove fluid from the bag. An example of this type of dispenser is given in Canadian Patent No. 1,206,923--Rutter (1986). In this dispenser, a semicircle is cut in the diaphragm, leaving a flap that is moved out of the dispenser fluid passage by fluid flow through it from the bag.\nAlthough the valves provided as part of either of these two types of dispensers are quite adequate for closing off fluid flow once use of the bag is commenced, the simple structure of such valves is usually not relied upon to prevent fluid loss during extended storage and shipment. That is why the bag is otherwise sealed during such storage and shipment, this seal being broken when the valve of the dispenser is first operated to an open position by the end user of the container.\nIt is an object of the present invention to provide a dispenser of the type having a sealing diaphragm therein with an improved structure for breaking the seal and allowing free fluid flow through the dispenser after the seal is broken."}
{"text": "Hitherto, there have been proposed and utilized many types of binding devices for filing papers and the like, which as a binding volume increases generally become unstable for binding operation, resulting in removal of bound papers. To file a great volume of papers or the like, therefore, special contrivances have been required for such binding devices.\nA binding device of the aforementioned type has been disclosed in the Japapese Laid-open Patent Publication No. 48-58915. Namely, according to one embodiment of this conventional binder particularly as shown in FIG. 11, the binding device comprise a fixture 10 and a cover plate 12 both made of synthetic resins, said cover plate 12 being provided with a slide opening 14 for slidably supporting a slide piece member 18 with engagement recesses 16, 16 so that the slide piece member 18 upon movement engages and disengages with a recess of a binding rod 20 for fixing and releasing the cover plate 12 from the fixture 10.\nFurthermore, according to another example as shown in FIG. 12, a slide plate 22 is integrally formed with a cover plate 24 and at the boundary portion between the slide plate 22 and the cover plate 24 is formed a hinge slit 26 for conveniently releasing the cover plate 24. The slide plate 22 is further provided with a slide slit 28 providing therein a fixing plate 30 which includes binding rods 34, 34 having engagement recesses 32, 32, while the cover plate 24 is provided at its external end with engagement piece members 36, 36 for engagement with the recess 32 of the binding rod 34 and the slide plate 22 and the cover plate 24 are moved together toward the fixing plate 30 to engage the binding rods 34, 34 with the cover plate 24. The fixture is further connected with the cover plate by means of a spring and under the spring function the cover plate may be enagaged with or removed from the binding rods.\nIn other words, according to such a conventional binding device the binding rods 20, 20 or 34, 34 for binding papers are firmly supported and engaged within a solid frame body comprising a fixture 10, a cover plate 12 and a slide piece 18 or a slide plate 22, a fixing plate 30, a cover plate 24 and engagement pieces 36, so that the filed papers, even if in plenty, may neither be unstable nor released.\nThe binding devices as hereinbefore described, however, usually consist of a plural member of components in combination, resulting in their complicated manufacturing process with a high cost and are further troublesome in filing operations.\nTherefore, an object of the invention is to provide a binding device which may be integrally molded and may provide a convenient binding operation."}
{"text": "Presently, electronic messaging constitutes a major form of interpersonal communications, complimentary to, and, in some respects, replacing, conventional voice-based communications. Electronic messaging includes traditional electronic mail (e-mail) and has grown to encompass scheduling, tasking, contact and project management, and an increasing number of automated workgroup activities. Electronic messaging also includes the exchange of electronic documents and multimedia content, often included as attachments. And, unlike voice mail, electronic messaging can easily be communicated to an audience ranging from a single user, a workgroup, a corporation, or even the world at large, through pre-defined message address lists.\nThe basic electronic messaging architecture includes a message exchange server communicating with a plurality of individual subscribers or clients. The message exchange server acts as an electronic message custodian, which maintains, receives and distributes electronic messages from the clients using one or more message databases. Individual electronic messaging information is kept in message stores, referred to as folders or archives, identified by user account within the message databases. Generally, by policy, a corporation will archive the message databases as historical data storing during routine backup procedures.\nThe information contained in archived electronic messages can provide a potentially useful chronology of historically significant events. For instance, message conversation threads present a running dialogue which can chronicle the decision making processes undertaken by individuals during the execution of their corporate responsibilities. As well, individual message store archives can corroborate the receipt and acknowledgment of certain corporate communications both locally and in distributed locations. And the archived electronic message databases create useful audit trails for tracing information flow.\nConsequently, fact seekers are increasingly turning to archived electronic message stores to locate crucial information and to gain insight into individual motivations and behaviors. In particular, electronic message stores are now almost routinely produced during the discovery phase of litigation to obtain evidence and materials useful to the litigants and the court. Discovery involves document review during which all relevant materials are read and analyzed. The document review process is time consuming and expensive, as each document must ultimately be manually read. Pre-analyzing documents to remove duplicative information can save significant time and expense by paring down the review field, particularly when dealing with the large number of individual messages stored in each of the archived electronic messages stores for a community of users.\nTypically, electronic messages maintained in archived electronic message stores are physically stored as data objects containing text or other content. Many of these objects are duplicates, at least in part, of other objects in the message store for the same user or for other users. For example, electronic messages are often duplicated through inclusion in a reply or forwarded message, or as an attachment. A chain of such recursively-included messages constitutes a conversation “thread.” In addition, broadcasting, multitasking and bulk electronic message “mailings” cause message duplication across any number of individual electronic messaging accounts.\nAlthough the goal of document pre-analysis is to pare down the size of the review field, the simplistic removal of wholly exact duplicate messages provides only a partial solution. On average, exactly duplicated messages constitute a small proportion of duplicated material. A much larger proportion of duplicated electronic messages are part of conversation threads that contain embedded information generated through a reply, forwarding, or attachment. The message containing the longest conversation thread is often the most pertinent message since each of the earlier messages is carried forward within the message itself. The messages comprising a conversation thread are “near” exact duplicate messages, which can also be of interest in showing temporal and substantive relationships, as well as revealing potentially duplicated information.\nIn the prior art, electronic messaging applications provide limited tools for processing electronic messages. Electronic messaging clients, such as the Outlook product, licensed by Microsoft Corporation, Redmond, Wash., or the cc:mail product, licensed by Lotus Corporation, Cambridge, Mass., provide rudimentary facilities for sorting and grouping stored messages based on literal data occurring in each message, such as sender, recipient, subject, send date and so forth. Attachments are generally treated as separate objects and are not factored into sorting and grouping operations. However, these facilities are limited to processing only those messages stored in a single user account and are unable to handle multiple electronic message stores maintained by different message custodians. In addition, the systems only provide partial sorting and grouping capabilities and do not provide for culling out message with duplicate attachments.\nTherefore, there is a need for an approach to processing electronic messages maintained in multiple message stores for document pre-analysis. Preferably, such an approach would identify messages duplicative both in literal content, as well as with respect to attachments, independent of source, and would “grade” the electronic messages into categories that include unique, exact duplicate, and near duplicate messages, as well as determine conversation thread length.\nThere is a further need for an approach to identifying unique messages and related duplicate and near duplicate messages maintained in multiple message stores. Preferably, such an approach would include an ability to separate unique messages and to later reaggregate selected unique messages with their related duplicate and near duplicate messages as necessary.\nThere is a further need for an approach to processing electronic messages generated by Messaging Application Programming Interface (MAPI)-compliant applications."}
{"text": "There are technical documents filed by the present applicant, which disclose techniques for blow-molding a cylinder-shaped molten resin into a skin-covered panel (see, for example, Patent Document No. 1: Japanese Laid-Open Patent Publication No. 10-235720).\nIn Patent Document No. 1, by means of the blowing pressure during blow molding, a skin material is heat-welded to an outer surface of one wall portion, and the other wall portion is formed with inner ribs protruding to contact an inner surface of the one wall portion, thereby obtaining the skin-covered panel.\nHowever, where the cylinder-shaped molten resin is used as in Patent Document No. 1, the molten resin is stretched in areas where the inner ribs are formed, and it is therefore necessary to increase the wall thickness of the cylinder-shaped molten resin in order to prevent formation of pinholes. This results in an increased weight of a final-molded-product skin-covered panel. Therefore, where the cylinder-shaped molten resin is used as in Patent Document No. 1, it may be difficult to reduce the weight and the thickness of the final-molded-product resin panel.\nNormally, the thickness (meaning “wall thickness”) of the cylinder-shaped molten resin (meaning “cylindrical parison”) is uniform. When split molds are clamped together, the pressing force for pressing the cylindrical parison against the split molds will be uniform across the entire surface of the cylindrical parison. Therefore, the parison pressed against one split mold for forming depressed portions such as inner ribs is stretched because of the blow ratio in accordance with the depressed portions, thereby resulting in locally thinned portions. Since no depressed portions are formed on the other split mold, no thinned portions will be produced. As a result, the thickness of the cylindrical parison needs to be set taking into account the thinned portions to be produced on the one split mold side on which the depressed portions are to be formed, thereby resulting in an unnecessary thickness on the other split mold side. Therefore, where a cylindrical parison having a necessarily uniform wall thickness is used, the wall thickness of the wall portion where no depressed portions are formed is greater than the wall thickness of the wall portion where depressed portions are formed after blow molding, and it is therefore difficult to reduce the weight and the thickness of the final-molded-product resin panel.\nNote that Patent Document No. 2 (Japanese Laid-Open Patent Publication No. 2010-155583) discloses a technique for obtaining a luggage board having a good mechanical strength property and rigidity even when the wall thickness of a resin panel of the luggage board is reduced for weight reduction.\nIn Patent Document No. 2, a rear wall is formed with a plurality of inner ribs integrally welded to a rear surface of a front wall, and a plurality of circular indentation-shaped round ribs integrally welded to the rear surface of the front wall are formed between the inner ribs of the rear wall. This makes it possible to obtain a required mechanical strength property and rigidity even when the average wall thickness of a resin wall of the luggage board is reduced to 1.5-2.5 mm for weight reduction. The average wall thickness of the resin wall of the luggage board is the average value between the wall thickness of the front wall and the wall thickness of the rear wall of the luggage board."}
{"text": "Tremolo or vibrato devices are well-known in the prior art. Tremolo devices are used in connection with guitars and similar musical instruments. Such devices serve to allow the player of a musical instrument to alter the tension of the strings of the musical instrument, thereby changing the pitch of the tone or sound produced by the strings.\nExamples of prior art tremolo devices may be found in various patents, such as, Fender U.S. Pat. No. 2,741,146; Rose U.S. Pat. No. 4,171,661; Rose U.S. Pat. No. 4,497,236 and Storey U.S. Pat. No. 4,457,201. In each of these devices a lever having a single arm is manipulated in a vertical direction in order to operate the device and alter the tension on the strings and raise the pitch of the notes produced by the instrument. For many users, the configuration of the current lever is less than desirable."}
{"text": "The Principles of the Tran-energy machines US patent US35,654 (WO 2009/112877) and PCT number PCT/IB200800796, which is incorporated herein by reference, describes the principles in the recovery of electric energy derived from the energizing of primary winding after producing the desired work, and harvesting the electric energy from secondary winding caused by increasing and decreasing in magnetic flux of primary winding during energizing and de-energizing.\nEPO application 14075049.8-1809, which is incorporated herein, describe a method and circuit for capturing electric energy utilized in energizing coils with capacitors and utilizing the capacitors as part of a voltage source to energize the coils.\nThe method that utilizes a time-varying current to build up a changing magnetic field in primary coil of a transformer. The changing magnetic field can:\n1—interact with another magnetic field to perform desired work(s);\n2—induce an inductive voltage across the secondary coil of winding on the main magnetic path;\n3—deflect the magnetic flux of permanent magnet in the main magnetic path, consequently magnetic flux of permanent magnets springs over air gaps to a subordinate (sub) magnetic path (less reluctance), therefore sub-inductive voltages appear across the sub-secondary windings, wrapped around the main-, sub-magnetic paths.\nHigh-potential electric charge from voltage-source, after releasing potential energy through a primary coil of winding, can be stored on capacitor, after that, the capacitor is in series with the voltage source, to supply electric charge with high-potential energy to the primary coil of winding. The same method of recovering the utilized electric charge can be applied for all instruments of the device.\nIn R-L or R-C or R-L-C electric circuits, at DC voltage-source, when switching ON respectively OFF, the current does not reach its desired amperes immediately, respectively zero ampere, but needs a short duration of time to change its values.\nChanging current causes changing magnetic flux. This is the basic for DC pulse transformer, moreover, with no voltage-drop full-wave bridge rectifier circuits at primary, secondary and sub-secondary coils of windings, electric energy can be regained, when magnetic flux in primary coil of winding decreases. Simultaneously more electric energy can be harvested at secondary coil of winding, and at sub-secondary coils of winding, wrapped around sub-magnetic paths.\nWhen charging or discharging a voltage-source that consists of capacitors in series or capacitor and battery in series, the same quantity of electric charge will flow through the voltage-source.\nHigh-potential electric charge can be formed and stored in low-leakage, low capacitance and high-voltage capacitor, not only by mean of a suitable capacitance, but also by Voltage-Doubler circuits, if needed. The theoretic operations of Voltage-Doubler and the well-known full-wave bridge rectifier circuit are not described here.\nHigh potential electric changes, stored on appropriate capacitors of primary, secondary, and sub-secondary coils of winding, can be used to recharge the selected voltage source by three ways:                1—Directly connected with the positive pole of the voltage source or via an electronic switch to control the following order of recharging;        2—Through a primary winding of a transformer, via an electronic switch, low-side of the transformer is connected with the positive pole of the voltage source. By this way, during the short time of transferring, the high-potential electric charge from capacitor can recharge the voltage source at the same time with inducing an inductive voltage across the secondary winding. Electric energy from the inductive voltage can be stored in appropriate capacitor for the next transferring.        3—The capacitor with high-potential electric charge is in series with an another voltage source that voltage of the one voltage source is equal with the selected voltage source. The in-series voltage source is connected with a primary winding of a transformer-containing permanent magnets via an electronic switch at high-side, low-side of the primary winding of the transformer is via an electronic switch, connected with the positive pole of the selected voltage source. When high-side and low-side switches are ON in a predetermined time, a desired time-varying current flows through the primary winding to recharge the voltage source, at the same time, it induces inductive voltages across all windings, related with the changing magnetic field of the primary winding. When high-side and low-side switches are OFF, electric energy in all windings flows to the appropriate capacitors, stored there and ready for the next transferring. The no voltage-drop full-wave bridge rectifier circuits or the voltage-doubler circuits are not described here.        \nWith two electronic switches, preferably without intrinsic body diode, at high-side and low-side, in OFF position, the primary coil of winding becomes an isolated voltage source, when regaining electric energy from collapsed magnetic field or the coming and going of magnetic field from permanent magnets, related with the primary winding.\nDirections of magnetic flux of permanent magnets can be deflected by interaction with changing magnetic flux of primary winding. Owing to this, more inductive voltages appear across the coils of winding, wrapped around the main magnetic core, the sub-magnetic cores. Therefore, more electric energy can be harvested.\nSoft magnetic materials with high magnetic permeability, low coercive force, can be used as magnetic wires. A multi-strand magnetic wire, combined of soft magnetic wire and for example, copper magnetic wires, or specially, magnetic wires, made of soft-magnetic material clad copper wire can be used to reduce the electric resistance and increase the strength of magnetic fields. Further drawing and description about these coils of winding, shall be done in the next application for patent of the Tran motors\\generators.\nFor electronic switches, controllable switching ON\\OFF in following order is essential to make very high-efficiency Tran-generators."}
{"text": "The headrest on an automobile seat is often made adjustable as to height. One common way in which the headrest is mounted on the seat in a way that permits height adjustment is for the headrest to be provided with two downwardly extending pegs, and the pegs engage sockets secured into the frame of the seat. Detent means are usually provided which interact between the pegs and the sockets, whereby the headrest may be set, by the occupant of the vehicle, at one of a number of pre-set heights.\nThe socket in which the peg is received comprises a metal tube. A plastic liner may be provided in the tube, to act as a bearing material for the peg. In the conventional system, the tube is welded to a bracket, and the bracket is welded to a frame piece of the seat. This manner of attachment, though secure enough (because it has to be secure by regulation), unfortunately is expensive as to the labour time and the materials needed to make it that secure. Any securement system that involves welding tends to be labour-intensive and therefore expensive, besides being difficult to inspect and test. A welded system generally has to be over-engineered.\nAlso, welding does not lead to high accuracy. The need for accuracy of placement of the headrest on the seat is not high, although the accuracy of the spacing of the pegs and their sockets cannot be too far out; the conventional welded-on system is just about at the limit for accuracy for welding, which means that, when welding is used as the basis of the attachment method, skilled care has to be taken, which in turn does nothing to ease the cost problem.\nThe invention is aimed at providing a manner of securing a headrest support tube into a seat frame, in a manner that eases some of the compromises that have had to be resorted to in the conventional systems.\nTypically, the operations carried out in a conventional seat manufactory include welding, bending of frame pieces, securing components together, and assembly, all of which tend to have a higher labour content. It is an aim of the headrest support system as described herein, to minimise the labour content of the task of attaching the support tubes to the seat frame piece.\nSupport tubes for headrests are conventionally attached to the seat frame piece by welding a bracket onto the frame piece, and then welding the tube to the bracket. Sometimes, the tube is pressed into holes in the welded-on bracket; but welding is nearly always resorted to, to assure that the tube remains in position on the bracket. Of course, the tubes can be attached securely enough, but the conventional costs of ensuring that security are high."}
{"text": "Modern large commercial buildings, such as factories, hotels, office buildings and hospitals, frequently use large and complex heating, ventilating and air conditioning (HVAC) equipment.\nIt is known to equip commercial buildings with variable air volume systems, which are capable of meeting the entire cooling and heating requirements of the building. Within the building, there are likely to be located a number of terminal units in different zones throughout the building, each connected via duct work to a central air supply. Such terminal units are sized to meet the conditions of the space which each serves, but, as a result, multiple offices, rooms or compartments within the structure are necessarily supplied with heating and cooling air by one terminal unit.\nThe end result of this type of design is that individual rooms or compartments within a structure are forced to share a common heating and cooling environment. While this may represent nothing more than a minor inconvenience for many building occupants in most cases, it presents particular difficulties in some specific environments, for example, hospital operating rooms.\nPrecise control of temperature and humidity in hospital operating rooms is important. Such rooms are frequently equipped with a number of machines which generate substantial heat. Further, the rooms will be populated with a varying number of workers during a typical operative procedure. Further, operating rooms must be regularly reconfigured for different procedures, meaning that the equipment and personnel contained within the room will vary substantially from day to day.\nUnder these circumstances, it is extremely difficult to maintain desired, consistent temperature and humidity levels in specific areas within buildings where centralized heating, ventilating and air conditioning systems are in use.\nWhile it is known to install modular heating, ventilating and air conditioning systems in individual rooms and compartments, many such devices are inefficient, cumbersome to install, and take up substantial space in the room in which they are installed. Further, such stand alone units are not centrally located within the rooms or compartments which they are designed to service, resulting in an imbalance in temperature and humidity in different areas of the same room or compartment. Further, such self-contained units often recirculate, rather than vent room air. Such units sometimes are in conflict in operation with the building central heating, ventilating and air conditioning system, resulting in energy inefficiencies when a local modular unit attempts to heat the air within a particular room or compartment at the same time as the centralized heating, ventilating and air conditioning system is attempting to cool the very same space.\nIt is desirable, therefore, to implement a modular heating, cooling and humidifying system which works in concert with the centralized heating, ventilating and air conditioning system of a larger structure, and which can be placed within the air ducting system of an existing structure, allowing individual temperature and humidity control in a single compartment or room, while at the same time not occupying physical space within the room or compartment, and further operating in symbiosis with the central heating, ventilating and air conditioning system of the structure."}
{"text": "1. Field of the Invention\nThis invention relates in general to accessing a database or executing a command over the Internet, and in particular, to a method and system for enhancing the performance of interpreted Web applications.\n2. Description of Related Art\nWith the fast growing popularity of the Internet and the World Wide Web (also known as xe2x80x9cWWWxe2x80x9d or the xe2x80x9cWebxe2x80x9d), there is also a fast growing demand for Web applications. For example, U.S. Pat. No. 5,737,592 describes a method and apparatus for executing SQL (Structured Query Language) queries in a relational database management system (RDBMS) via the Web. Web users can request information from the RDBMS via HTML (HyperText Markup Language) input forms, which request is then used by an interpreted macro language file to create an SQL statement for execution by the RDBMS. The results output by the RDBMS are themselves transformed using the macro language file into HTML format for presentation to the Web user.\nAn application that is interpreted when executed offers a high degree of flexibility to developers and installers of the application because any changes that are made to the application take effect immediately. An interpretive execution environment also offers a high degree of flexibility to users, because users can provide input that alters the execution of the application or even make dynamic changes to the application itself at the time they request execution of the application.\nHowever, the performance of interpreted applications is constrained by the parsing costs that are typically incurred for each execution of the application. These costs can be considerable and may limit the usefulness of the applications within production environments, and in particular, within production environments where customer acceptance of the applications is dependent on acceptable response time. Thus, the value of the flexibility described above is lost.\nThere is a need in the art for techniques that improve the performance of interpreted Web applications without sacrificing the flexibility that this style of application makes available to application programmers, installers, and users of the application.\nTo overcome the limitations in the prior art described above, and to overcome other limitations that will become apparent upon reading and understanding the present specification, the present invention discloses a system, method, and article of manufacture for caching a pre-processed version of an interpreted application at the time the interpreted application is first interpreted and executed by a computer, thereby improving response time for subsequent executions of the interpreted application. The pre-processed version of the interpreted application captures all execution paths through the interpreted application, despite the fact that the first execution of the interpreted application may only exercise a subset of the possible execution paths therein. Modifications to the interpreted application are dynamically detected, so that the pre-processed version can be replaced in the cache as needed.\nReferring now to the drawings in which like reference numbers represent corresponding parts throughout:\nFIG. 1 schematically illustrates the hardware and software environment of the preferred embodiment of the present invention;\nFIG. 2 schematically illustrates the structure of a macro language file according to the preferred embodiment of the present invention;\nFIG. 3 schematically illustrates the structure of a hash table for macro language files according to the preferred embodiment of the present invention;\nFIG. 4 schematically illustrates the structure of each entry of the hash table according to the preferred embodiment of the present invention;\nFIG. 5 schematically illustrates the structure of a mutex array for the hash table according to the preferred embodiment of the present invention; and\nFIGS. 6, 7, and 8 are flowcharts illustrating the logic performed by the Net.Data gateway in interpreting macro language files according to the preferred embodiment of the present invention."}
{"text": "1. Field of the Invention\nThe present invention relates to self-drilling anchors. In particular, it relates to a self-drilling anchor assembly and system for suspending a member from a hollow wall.\n2. Description of the Related Art\nConventional hollow wall anchors are well known in the art. These anchors have evolved along a common conception in the design of anchoring members to a hollow wall. One such design is the toggle bolt anchor, exemplified by the TOGGLER® anchor, which when positioned, expands behind the inner surface of a wall, to spread the holding load against a greater surface of the wall. In this manner, the toggle bolt anchor provides an ability to support a heavy load. However, this type of anchor requires pre-drilled holes for installation and often damage to the wall, upon removal. Self-drilling anchors are also well known. These anchors generally include a plastic tapered body having external helical threads with a contiguous driving end for taping into a hollow wall without the need for a predrilled aperture in the wall. Up until now, a few such examples of self drilling anchors exist.\nOne such example, U.S. Pat. No. 6,196,780 to Wakai et al., discloses a threaded anchor having a hollow shank. The shank has a head at one end, an integral blade portion at the other end, and is formed with external threads on the outer periphery. Cuts are formed in a peripheral wall for the shank to define a bendable portion. The bendable portion is adapted to bend when pushed out by a mounting screw driven into the hollow shank. The blade portion is semicylindrical with side edges defining cutting edges. Tips for inclining the mounting screw are formed on the inner surfaces of the blade portion and the hollow shaft, respectively. When a mounting screw is driven through the article into the anchor to fix an article to the board, the mounting screw is inclined relative to the anchor, so that the bendable portion is pushed out, thus preventing the anchor form rotating together with the screw.\nAnother example, U.S. Pat. No. 7,144,212 to Kaye et al., discloses a self-drilling hollow wall anchor. The anchor has a self drilling threaded installation body adapted to be broken, after positioning, by inserting a screw into expandable anchoring elements. The anchoring elements are designed to engage internally formed cams and the distal surface of the wall into which the anchor is deployed. The anchoring elements are made of a plastic which is designed as a compromise such that it is sufficiently ridged in order to effect a threaded insertion and also to be broken apart. The anchor also includes anti-rotation elements to retard over-tightening and anchor removal prior to screw removal.\nWhile the foregoing examples offer some utility, a major disadvantage in each lies in the fact that, while they do provide for a self drilling anchor, they do not provide the high load bearing capacity of a toggle bolt anchor, and their one piece plastic design is inherently compromised in providing a drilling tip strength but are made of sufficient flexibility so as to allow for expansion of the internal anchoring members. The foregoing examples are also difficult to remove and are limited in application for use in the anchoring of members which can be suspended through the head of a screw. Thus, it is desirable to provide a self-drilling anchor which is designed to provide the high load bearing capacity in conjunction with a hardened pin-point drilling tip for accurate placement, but which is also easily removed without damage to the wall and which is capable of suspending members other than through the head of a screw. The present invention satisfies these needs."}
{"text": "Malaria has a tremendous impact on human health, killing millions annually and the disease is a major impediment for social and economic development of nations in malaria-endemic areas, particularly in sub-Saharan Africa (Sachs & Malaney (2002) Nature 415:680-85). Malaria is a mosquito-borne disease that is transmitted by inoculation of the Plasmodium parasite sporozoite stage. Sporozoites invade hepatocytes (Kappe et al. (2003) Trends Parasitol. 19:135-43), transform into liver stages, and subsequent liver stage development ultimately results in release of pathogenic merozoites that initiate the blood stage cycle of infection (Shortt & Garnham (1948) Nature 161:126).\nProtection against blood stage malaria can be passively transferred by antibodies. Effectiveness of passive transfer of IgG between adults and children living in different geographic regions indicate that the antigens that are targeted by antibodies that protect against blood stage malaria are conserved (see Duffy et al. (2005) Vaccine 23 (17-18):2235-42). The best evidence that naturally occurring immunity against blood stage malaria targets the IRBC surface comes from studies of pregnancy malaria. In areas of stable malaria transmission, adults enjoy immunity that limits parasitemia and prevents disease. Women become more susceptible during pregnancy, and previously this was ascribed to pregnancy-related immunomodulation that develops to prevent rejection of the fetal allograft. However, susceptibility is greatest in primigravid women and diminishes over successive pregnancies, suggesting an acquired immune response to an antigenically distinct microbe. Placental isolates of P. falciparum commonly bind to chondroitin sulfate A (CSA) expressed on the surface of the syncytiotrophoblast (the cells lining the maternal vascular space in the placenta), and this phenotype is uncommon among isolates obtained from non-pregnant individuals (Fried & Duffy (1996) Science 272:1502-4). The distinct binding phenotype renders pregnant women naïve to this parasite population during their first pregnancy (primigravidas). Women with multiple pregnancies (multigravidas) residing in areas of stable malaria transmission develop antibodies that inhibit parasite adhesion to CSA (Fried et al. (1998) Nature 395:851-2). These anti-adhesion antibodies are associated with a reduced mass of parasites in the placenta, and substantial improvements in fetal development (Duffy & Fried (2003) Infect. Immun. 71:6620-3). Because CSA-binding parasites cross-react with sera donated by mothers throughout Asia and Africa, the antigen targeted by protective antibodies is presumed to have conserved features or limited diversity.\nThere is a need in the art for vaccines that protect against malaria infection and disease. There is also a need in the art for diagnostic markers for malaria. The present invention addresses these needs and others."}
{"text": "In recent years, development of technology for wirelessly transmitting electric power (electric energy) without using power supply cords and the like has become popular. Among methods for wirelessly transmitting electric power, the technique called a magnetic resonance method is gaining particular attention. The magnetic resonance method was proposed by a research group at the Massachusetts Institute of Technology in 2007. The related technology is disclosed in Patent Document 1 (JP2009-501510A), for example.\nIn a wireless electric power transmission system of the magnetic resonance method, a resonance frequency of a power transmission antenna is equal to a resonance frequency of a power receiving antenna; the antennas with a high Q-value (100 or more) are used. Therefore, from the power transmission antenna to the power receiving antenna, energy can be efficiently transmitted. One of major features is that the power transmission distance can be from several tens of centimeters to several meters.\nStudies have been under way on how to apply the above-described magnetic resonance-type wireless electric power transmission system to a process of charging vehicle-mounted batteries, such as those of an electric vehicle (EV) or a hybrid electric vehicle (HEV). The use of such a wireless system for the above vehicles would eliminate the need to handle power supply connectors, power supply lines, and the like for the supply of electric power to the vehicles.\nFor example, what is disclosed in Patent Document 2 (JP2010-68657A) is one antenna that is mounted in a bottom surface portion of a moving object such as an electric vehicle, with another antenna provided on the ground to wirelessly transmit electric power to charge a battery of the electric vehicle with the transmitted electric power.\nSeveral proposals have been made on a specific configuration of an antenna used in the magnetic resonance-type wireless electric power transmission system. For example, as for the configuration of communication coils that are provided in a power feeding circuit and a power receiving circuit in a wireless electric power transmission device that wirelessly transmits electric power from the power feeding circuit to the power receiving circuit, what is disclosed in Patent Document 3 (JP2010-73976A) is the configuration of the communication coils of the wireless electric power transmission device characterized by including: a printed board made of a material with a relative permittivity of more than 1; a primary coil that is provided on a first layer of the printed board and is formed as a conductive pattern making at least one loop; and a resonance coil that is provided on a second layer of the printed board and is formed as a spiral conductive pattern.    Patent Document 1: JP2009-501510A    Patent Document 2: JP2010-68657A    Patent Document 3: JP2010-73976A"}
{"text": "The present invention relates to slings for firearms and the like, and particularly to a sling for carrying a weapon in a generally horizontal position of readiness.\nIn order to be able to use shoulder weapons such as hunting rifles and the like as quickly as possible, yet have freedom to use their hands, many hunters would prefer to carry such weapons in a forwardly or laterally directed position with the weapon upright and its barrel or main longitudinal axis approximately horizontal. This position is known as a horizontal ready position.\nRifles, particularly military rifles, have long been equipped with slings permitting them to be carried without encumbering the carrier's hands. Such slings are also useful in bracing such weapons during use, to achieve a better aim. Slings have conventionally been located attached to the bottom of a rifle or other shoulder-fired weapon, at least partly because location of a sling mounting at the top of the weapon might interfere with the use of sights mounted low on the top of the weapon. Conventional military type slings, however, do not permit a weapon to be carried suspended by the sling in the preferred horizontal ready position.\nMoving the point of attachment of a sling to the top of a weapon so that it might be carried in a horizontal ready position has previously resulted in the sling interfering with the ability to raise the weapon quickly to the proper position against the user's shoulder. For example, when a sling is attached at the top of the rear end of the butt stock, raising the weapon to a horizontal position with the butt end of the stock against the user's shoulder might result in the sling being caught between the stock and the user's shoulder or lying across the top of the cheek rest area of the stock, where it would adversely affect the user's ability to aim the weapon quickly and accurately. Since the ability to carry a weapon in the horizontal ready position can provide a definite advantage in the ability to shoot quickly, either in infantry warfare or in hunting, it is desirable to provide a sling which can be used to carry a rifle in the horizontal ready position, but which does not interfere with the use of the weapon's sights.\nA previously attempted answer to these needs was an adaptor kit for attaching the front of the sling strap to the front sight assembly of a military rifle, with the rear end of the sling fastened as a noose extending through the conventional sling swivel and around the butt stock. This, however, left the rear of the sling in an undesirable location during use of the weapon, and the front attachment was likely to be noisy.\nParticularly with hunting rifles, it is not uncommon to have stocks made of expensive wood, beautifully finished. Owners of such rifles are, naturally enough, reluctant to attach a sling which might damage the finish of the stock. It is therefore desirable to provide a sling which can be used to carry a hunting rifle in the horizontal ready position without marring the finish of its stock.\nFor both hunting and military weapons it is desirable to avoid any noisy attachments.\nBrokus U.S. Pat. No. 3,606,109 discloses a gun sling including a socket loosely fitted to the underside of the rear of the stock, to facilitate carrying a gun in a horizontal ready position. The socket taught by Brokus, however, is undesirably bulky and apparently does not remain in place unless tension is maintained in the sling.\nMoomaw U.S. Pat. No. 2,915,233 teaches a gun sling including a stock-supporting loop or socket which enables a gun to be carried in a horizontal ready position. Raising the gun from a horizontal ready position to a normal firing position with its stock against the user's shoulder, however, completely disconnects the rear end of the sling from the stock, requiring it to be replaced on the stock before the sling can be used subsequently to carry the gun in any position whatsoever.\nJohnson U.S. Pat. No. 4,555,051 discloses a sling assembly of the desired type which is particularly adapted to certain military weapons equipped with front sight assemblies that are large enough to be used as a point of attachment for a front end of the sling strap. The sling assembly disclosed, however, is not particularly well adapted to weapons not equipped with such a high front sight assembly.\nWhat is needed, therefore, is a sling assembly which can be used to carry a weapon such as a hunting rifle quietly, in a position of readiness in which the weapon is upright with its longitudinal main axis directed generally horizontally, in which the sling assembly does not interfere with the ability to raise the butt end of the weapon's stock properly to the user's shoulder, nor with the ability to use the installed sights of the weapon, and in which the sling assembly is not likely to mar the finish of the weapon or its stock."}
{"text": "The IPD method comprises placing a cathode formed of a coating material into a vacuum environment within a vacuum chamber, providing a workpiece having a surface within the vacuum environment, supplying a current to the cathode to form a cathodic arc upon a cathode surface resulting in erosion or evaporation of coating material from the cathode surface, and depositing the coating material from the cathode upon the workpiece surface. The surface of the workpiece may be referred to more generally as a substrate surface. The workpiece may be a turbine blade part or component in need of coating or repair, but the method and apparatus can be used to deposit coating materials upon any substrate surface that can be subject to an IPD method.\nA cathodic arc is formed upon a surface of a cathode when a cathode and an anode are placed in a vacuum environment, and a sufficient voltage is applied to the cathode to create a sufficient potential between the cathode and the anode to form a discharge on a surface of the cathode. This discharge creates a highly ionized metal plasma from the erosion of the cathode. The arc is formed for a short period of time, extinguishes and reforms rapidly, giving the appearance that the arc is moving continuously over the surface of the cathode.\nThe erosion of the cathode forms a cloud of coating material containing ions, charged particles, vapors and neutral droplets in the vacuum environment. This coating material is then available for deposition within the vacuum environment. The coating material may deposit upon cool surfaces in the vacuum environment by condensation, and also upon an anode surface by electrical attraction as well as condensation.\nA workpiece may be placed in the vacuum environment to form a coating upon a surface of the workpiece. The workpiece may be manipulated within the vacuum environment to produce a uniform coating over the entire surface of the workpiece, or the workpiece may be masked in order to produce a coating on a selected surface. The workpiece may be held at ground potential or a voltage bias may be applied to the workpiece to affect material deposition by increasing the attraction of ions and charged particles to the surface.\nThe anode may be a conductive surface of the vacuum chamber or a separate structure within the vacuum chamber. The function of the anode is to sink electrons from the cathode to sustain the electrical discharge.\nIn the past, cathodes have been formed of solid cylindrical stock having a flat end-face. These cathodes are typically cut to short lengths of about 4.4 cm with an outside diameter of about 6.3 cm. The plasma arc is formed at the flat end-face of the cathode open to the vacuum environment.\nMultiple cathodes are often positioned in a vacuum chamber around the workpiece to increase the amount of coating material available for coating upon the workpiece surface. These conventional cathodes operate on current of approximately 40 A to approximately 1200 A.\nA plasma deposition system has been disclosed by S. A. Mouboiadjian, E. N. Kablov and Ya. A. Pomeloov in their paper “Equipment, technology and protective coatings produced by ion bombardment deposition.” This system discloses a cylindrical cathode that has an arc formed on a circumferential surface. However, the structure requires a ring-shaped cooled anode and electromagnetic coil fixed coaxially around the cathode to effect arcing.\nCurrently, the cost of an IPD process is dominated by factors pertaining to the cathode. A cathode that allows for more coating material to be available to the coating process will have a significant impact on the process cost by allowing the process to be operated for greater periods of time without having to replace the cathodes. Additionally, by not requiring the cathode to be replaced as often, a more uniform and improved coating material may be deposited.\nOther features and advantages of the present invention will be apparent from the following more detailed description of a preferred embodiment, taken in conjunction with the accompanying drawings, which illustrate by way of example, the principles of the invention."}
{"text": "Fibre optic distributed acoustic sensors are known. Such sensors interrogate optical fibres with optical radiation and measure changes to the radiation resulting from acoustic waves affecting the optical fibre.\nU.S. Pat. No. 5,194,847 (Texas A&M Univ) describes interrogating a fibre with a repeated coherent pulse of radiation and detecting any radiation which is Raleigh backscattered within the fibre. The fibre is interrogated with a single pulse at a time and the amplitude of the backscattered radiation is analysed to detect any disturbance of the fibre by acoustic/pressure waves. This document teaches that a buried optical fibre can be used as a distributed acoustic sensor for perimeter monitoring purposes.\nUK Patent Application GB2,222,247 (Plessey) describes another distributed fibre optic sensor system in which changes in environmental parameters, such as sound waves, are sensed by transmitting pulses of light along an optical fibre. This document describes that two closely spaced pulses may be transmitted into the fibre, the first pulse having a different frequency to the second pulse. The backscatter from the pulses within the fibre can be detected and analysed at a carrier frequency equal to the frequency difference between the interrogating pulses. The signals received at a detector can be gated and processed to determine information representative of changes in environmental parameters affecting a desired section of the optical fibre.\nUK patent application GB2,442,745 (AT&T) describes distributed acoustic sensing using an optical fibre. This document again teaches use of pulse pairs, wherein the individual pulses of a pulse pair have different frequencies. The backscattered signal is analysed at a carrier frequency corresponding to the frequency difference between the pulses in the pulse pair. This document teaches applying a complex demodulation to the detected backscatter signal at the known frequency difference between the pulses in a pulse pair to provide in-phase (I) and quadrature (Q) signals for the carrier frequency. These are then converted to provide the phase and amplitude of the signal. The phase of successive samples from the same section of fibre is then monitored to determine any acoustic signals incident on that section of fibre.\nThis document (GB2,442,745) teaches that the frequency difference between pulses in a pulse pair should be related to the pulse width. The example is given of pulses of 20 m width and a frequency difference between pulses in a pair of at least 5 MHz.\nWhilst the technique described in GB2,442,745 is useful, in some instances the baseband structure inherent in such a fibre optic sensor, i.e. a random but systematic pattern in the detected backscatter, can mask or destroy the carrier signal and reduce the signal to noise ratio of the sensor. This baseband structure of the system arises partly from the random distribution of the scattering sites in the optical fibre, from thermal drift etc. and thus can not be eliminated. The effect of cross over of the measurement signal and baseband noise of the system can be mitigated by using higher carrier frequencies, for instance of the orders of hundreds of MHz. However use of such a high carrier frequency would require detector sample rates in excess of hundreds of MHz. This not only would require very fast components for the interrogator unit and greatly increase the amount of processing required but a much higher detector bandwidth time would also impact on the sensitivity of the detector."}
{"text": "American and international telecommunications service providers today purchase termination services from a large number of carriers around the world. As previously centralized national telecommunications networks have opened to increasing competition, the number of network routes available for the termination of both domestic and international traffic has increased worldwide. This trend has been especially apparent since the mid-1990s during which time the deregulation of parts of American and other telecommunications networks engendered a dramatic increase in the number and quality of independent telecommunications carriers servicing the American market.\nAdditionally, foreign telecommunications markets have been deregulated and telecommunications markets have transitioned from the single government-owned incumbent carriers to a collection of multiple separate carriers.\nThis increased carrier penetration of different world telecommunications markets has given American telecommunications service providers a greater number of options for the termination of international telecommunications traffic than was previously available, along with lower costs of termination to many international destinations, especially in the developing world. Unfortunately, however, a growing number of the less costly options now available to American service providers for the termination of international traffic frequently fail to deliver a quality of service comparable to more costly established international carriers. Thus American and other telecommunications service providers interested in the cost savings associated with smaller, independent international carriers must guard against the potentially damaging impact of poorly qualified secondary routes to customer perceptions of quality.\nExisting systems and methods to measure and manage quality of service are limited in several important ways. Some quality of service measurement and management methods equate network quality with a measure of Answer-Seize Ratio (ASR) either in part or in whole. However, experience has shown that small, independent carriers providing low-cost routes to the developing world may accidentally misrepresent or intentionally falsify release cause codes transmitted as part of the SS7 and ISDN telecommunications signaling protocols. For example, carrier substitution of a network congestion cause code in place of a user busy cause code may reclassify an unanswered call as a non-attempted call, thereby raising the measure of ASR for the carrier in question. Another limitation of quality of service measurements relates to ASR, PSC (percentage of short calls) and similar measures which express call quality as a ratio over total calls and can give inaccurate results over small samples. For example, a single caller repeatedly calling the same number to some low volume international destination may disproportionately raise or lower measures of ASR or PSC for the destination over a small sample of calls. As quality of service correction in realtime sometimes necessitates relatively small call samples, realtime quality correction may confound quality of service management systems predicated on ASR or PSC alone.\nThus, there exists a need to provide new and improved systems and methods that permit telecommunications service providers to measure and manage end-user perception of quality of service. A need has arisen to effectively address the needs of modem telecommunications service provides, in which such systems and methods must measure voice traffic quality of service to places in both the developed and developing worlds without reference to the special messaging parameters included in SS7, ISDN and other signaling protocols, and must correct for the skewing of repeated caller-to-caller dialing and must further afford comparison between different carriers terminating to the same destination."}
{"text": "For the purposes of operation and in order to store information, computer systems and mainframes have a memory arrangement comprising a multiplicity of memory modules which are arranged, for example, in memory cabinets and have semiconductor components grouped on the memory modules.\nSemiconductor components, for example, DRAMs (Dynamic Random Access Memories), are generally subject to extensive function tests early in the production process before final assembly to form a memory module. These function tests are used to identify faulty memory cells or faulty column lines or row lines or generally faulty circuit parts in the semiconductor components. To this end, data values are written to memory cells in a memory cell array in the semiconductor component and are then read out again in order to be compared with the prescribed data values. This makes it possible to test the semiconductor components under various operating conditions in order to guarantee fault-free operation of the memory chip.\nThe memory modules accommodate a test device, a “BIST” (Built-In Self-Test) unit, in each of the semiconductor components as part of the latter. The BIST unit integrated in the respective semiconductor component carries out the requisite electrical function tests before the semiconductor components are installed. The BIST unit has a BIST controller which, as a switching region in the semiconductor component, is in the form of an ASIC (Application-Specific Integrated Circuit). Commands in a test sequence which are issued by the BIST controller are forwarded to the semiconductor component, with the BIST controller monitoring and evaluating the execution of the commands. The data transmitted by the semiconductor component regarding its operating states are output, for example, to external test systems which make an appropriate evaluation on the basis of which it is possible to make a statement regarding whether and, if appropriate, which memory areas are not functioning as intended. When the tests are carried out successfully, the module is classified as functional and is used in the customer's target application.\nHowever, a meaningful test result can only be achieved when the semiconductor component is tested at the operating frequency that it has during normal operation. A fault in a semiconductor component is always associated with the target application, for example, a voltage supply or an input parameter for configuring the semiconductor component is not within the prescribed specification.\nHowever, today's test methods do not yet make it possible to simulate these or all characteristic operating modes of the application in order to test the semiconductor components in proximity to the application. It is thus not possible to make a statement regarding whether the semiconductor components tested during production will run through all of the operating modes, occurring during later application without any faults.\nIn the event of a fault occurring during normal operation at the customer's premises, the memory module has to be returned to the manufacturer for the purposes of analysis. Identification data which may have been programmed in, for example the chip ID, test data or adjustment parameters, can thus be used only for subsequent historical tracking, but not for user-specific adjustment during normal operation.\nAt such a point in time it would be desirable, for the purposes of evaluation and/or analysis, to use a test system which would make it possible to test and adjust the semiconductor components during normal operation. Today's available external test systems are connected to the semiconductor components on a memory module via the semiconductor component's standard interface, which is used for external data interchange, address interchange and/or command interchange during normal operation. In the test mode, the test system is able to generate the test commands required to test the memory module, such as control and address commands, commands for reading and storing data words and also a clock signal, and is able to initiate the electrical function test via the BIST unit, for example. In the case of semiconductor components which can be operated in parallel, however, this function test may usually be carried out only for all chips, that is, all of the semiconductor components arranged on the memory module are tested in parallel at the same time via the standard interface.\nThe desire for remote access monitoring and a fault analysis capability for computer systems, i.e., driving the computer systems via an external test system, has existed since computer systems were networked. One difficulty in checking individual hardware components in a computer system is based on the fact that respective redundancy must be available for faulty components.\nFor example, an apparatus for the remote monitoring of computer components in a computer system is known, but this apparatus explicitly requires a working memory in the computer system. In order to analyze a processor in a computer system, a remotely accessible integrated debug environment has been proposed, wherein a computer which is connected to the computer system via the Internet can analyze the processor in the event of a fault. Further, a method for remotely accessing a faulty booting computer is known, in which the computer, in the event of a failed starting attempt, has recourse to a simple E-BIOS code which connects the computer to a service computer via a LAN or an Internet connection and thus makes the computer accessible for remote access operations for the purpose of repair and/or diagnosis. A known method for testing an SDRAM in a computer system, including test modules integrated in the computer system, tests the memories before starting by appropriate test modes in order to boost or attenuate a possible fault mode.\nIn addition, there is known a semiconductor module having semiconductor components which are arranged on the semiconductor module and are connected to one another via a serial line and having an interface which accesses the semiconductor components via the serial line.\nHowever, remote access maintenance of memory modules during normal operation has not been disclosed. To this end, it would be necessary to operate the computer system normally and to carry out specific addressing operations during an application.\nA memory arrangement which makes it possible to drive the semiconductor components (arranged on the memory arrangement) during normal operation using remote access and to test and adjust them in proximity to the application without impairing the operation of the semiconductor components is desirable."}
{"text": "The present invention relates to a starting device for a discharge lamp such as a fluorescent lamp. More particularly, the invention relates to a discharge lamp starting device utilizing semiconductor switching elements.\nVarious starters utilizing semiconductor switching elements have been heretofore proposed, of which one example utilizing a nonlinear dielectric element and a thyristor is shown in FIG. 1. In this figure, reference numeral 1 designates a discharge lamp including filaments 101a and 101b at opposite ends of the lamp; reference numeral 2 designates an inductive stabilizer; reference numeral 3 a semiconductor switch composed of a reverse blocking triode thyristor 301, a trigger element 302 such as an SBS (Silicon Bidirectional Switch) or Diac, voltage dividing gate circuit resistors 303a and 303b, and a smoothing capacitor 304; reference numeral 4 designates a nonlinear dielectric element; reference numeral 5 a noise eliminating capacitor; and reference characters U and V power source terminals.\nIn the above-described starter, when an AC voltage e.sub.UV, having a waveform shown by a dotted line in FIG. 2A, is applied across the power source terminals U and V, at the beginning of the lamp starting period, the thyristor 301 turns on at a phase angle .theta..sub.1 in a positive half cycle of the power source voltage. At that time, a current flows through a path including the stabilizer 2, filament 101a, thyristor 301, and filament 101b, thereby preheating the filaments. This current will lag the source voltage due to the inductive effect of the stabilizer 2. When the current flowing through the thyristor falls below the holding current of the device, the thyristor 301 turns off at a phase angle .theta..sub.2, which occurs during the negative half cycle of the power source voltage. At the instant that the thyristor 301 is turned off, the voltage applied across the element 4 is zero. Since the power source voltage e.sub.UV is then in the proximity of the negative peak of the waveform, the element 4 is subsequently charged to the indicated polarity through the stabilizer 2.\nThe element 4 has a Q-V (stored charge vs. voltage) as shown in FIG. 3, wherein the stored charge becomes saturated at an applied saturation voltage E.sub.s. By appropriately selecting the Q-V characteristic of the element 4, the nonlinear region (where the charged voltage is less than the saturation voltage E.sub.s) is reached at a voltage less than the peak voltage of the power source. In this case, the charging current flowing into the element 4 is abruptly reduced at the instant the power source voltage exceeds the saturation voltage. Due to the inductive property of the stabilizer 2, the charged voltage of the element 4 then increases abruptly in the form of a pulse voltage V.sub.21, as shown in FIG. 2A which is substantially higher than the peak voltage of the power source. The pulse voltage V.sub.21 is applied across the discharge lamp 1. After the occurrence of this pulse voltage, the power source voltage e.sub.UV is applied across the lamp 1 until the thyristor 301 again turns off in the next cycle.\nThe above-described operation continues until the lamp 1 is started. That is, while the filaments 101a and 101b of the lamp 1 are being heated by the preheating current, the discharge of the lamp 1 is initiated by one of the positive pulse voltage V11 and the negative pulse voltage V21.\nOnce the lamp 1 has started, the lamp voltage is reduced below the power source voltage, thus keeping the thyristor 301 turned off. More specifically, although the lamp voltage instantaneously rises above the power source voltage, as shown by V.sub.12 and V.sub.22 in FIG. 2A, due to the charging effect of the element 4, the thyristor 301 cannot be turned on by voltages of the magnitude of V.sub.12 because of the smoothing effect of the capacitor 304.\nWhile the above-described starting device utilizing a nonlinear dielectric element and a thyristor is advantageous in its starting performance, simplicity of circuit construction, and low cost, nevertheless, the construction shown in FIG. 1 has following difficulties:\n(1) The power consumption of the circuit is higher than in a starter in which the starter is disconnected from the lamp circuit after the lamp has been started.\n(2) Charging and discharging currents flowing in and out of the element 4 create annoying vibration and noise due to piezoelectric effects.\nMore specifically, as shown in FIG. 2B, after the lamp 1 has been started (at phase angles .theta..sub.7 and .theta..sub.8, for instance), a discharge current i.sub.11 flows through the element 4 and the lamp 1. Due to the presence of the discharge current i.sub.11, the power consumption of the lamp 1 is considerably high in comparison with starter arrangements in which the starter is fully disconnected from the lamp circuit after the lamp has been started.\nFIG. 4 illustrates a second example of a conventional starting device which is intended to overcome the above-described difficulties. In this figure, reference numeral 30 designates a bidirectional semiconductor switch, reference numeral 7a designates a diode connected in parallel with the element 4, and reference numerals 6 and 7b designate a series-connected resistor and diode connected in parallel with the semiconductor switch 3 to provide a discharge circuit for the switch 30. The semiconductor switch 30 is composed of a bidirectional triode thyristor 305, a trigger element 302 such as an SBS, Diac or the like, resistors 303a and 303b, and a capacitor 304.\nThe operation of the conventional device will now be described with reference to FIG. 5A which shows a voltage waveform across the lamp 1.\nIn the beginning of the starting operation, the thyristor 305 is turned on at a phase angle .theta..sub.1 in a positive half cycle of the power source voltage e.sub.UV, at which time a preheating current flows through a path including the stabilizer 2, filament 101a, diode 7a, thyristor 305, and filament 101b. The thyristor 305 is turned off at a phase angle .theta..sub.2 in the following negative half cycle of the power source voltage e.sub.UV, at which time the preheating current is reduced to zero. The thyristor 305 is again turned on thereafter at a phase angle .theta..sub.3 by operation of the trigger element 302, thereby to create a charging current through the thyristor 305 and the element 4.\nBecause the element 4 has a nonlinear characteristic as in the case of the first example of the conventional device, when saturation is reached and the charging current of the element 4 drops at a phase angle .theta..sub.4 to a value lower than that required for maintaining the thyristor 305 in the conductive state, the thyristor 305 is again turned off and the power source voltage is continuously applied across the lamp 1 until a phase angle .theta..sub.6 in the positive half cycle, at which time the thyristor 305 is again turned on to pass the preheating current.\nThe purpose of the discharge resistor 6 and the diode 7b is as follows. When saturation is reached at the phase angle .theta..sub.3 causing an abrupt increase of the voltage across the element 4 to the maximum voltage V.sub.21, the current through the element 4 then flows through the resistor 6 and the diode 7b, thereby causing the voltage across the element 4 to substantially follow the lamp voltage. Since the voltage applied to the thyristor 305 is substantially equal to the difference between the charged voltage V.sub.21 of the element 4 and the power source voltage e.sub.UV, if there were no such discharge circuit, an extremely high voltage withstanding property would be required for the thyristor 305. Also, the diode 7b prevents charging of the element 4 during the time interval between phase angles .theta..sub.2 and .theta..sub.3, thus ensuring the generation of the high voltage pulse V.sub.21 by abruptly charging the element 4 starting from a zero potential.\nOnce the discharge lamp 1 has started, the lamp voltage is reduced below the power source voltage e.sub.UV, thus preventing turning on of the thyristor 305 and maintaining stable operation of the lamp 1. Furthermore, since most of the power source voltage is applied across the thyristor 305, the voltage applied to the element 4 is reduced to approximately zero. Hence the drawbacks hereinbefore described with respect to the first example of the conventional device that are caused by charging and discharging currents of the element 4 after the lamp has been started are eliminated.\nUnfortunately, however, as illustrated by FIG. 5B, the connection of the diode 7a in parallel with the element 4 prevents positive voltages from being applied across the element 4. This causes the dielectric polarization of the element 4 to be shifted in one direction only, and hence the hysteresis loop of rectangular shape as shown in FIG. 3 is deformed to such an extent that the desired nonlinear characteristic of the element 4 is substantially lost. In other words, the amplitude of the pulse generated at the phase angle .theta..sub.4 is reduced to such an extent that the starting of the discharge lamp is difficult.\nAn object of the present invention is thus to overcome the above-described drawbacks of the conventional discharge lamp starting devices.\nMore specifically, it is an object of the present invention to eliminate:\n(1) the high power consumption due to charging and discharging currents of the element 4 and the accompanying generation of vibration from the element 4; and\n(2) the reduction of negative-side pulse voltages from the element 4."}
{"text": "1. Field of the Invention\nThe present invention relates to a nonvolatile semiconductor memory, and in particular, to resistors in a nonvolatile semiconductor memory having a stack structure and a method of manufacturing the resistive elements.\n2. Description of the Related Art\nNonvolatile semiconductor memories use elements of a fixed resistance to generate internal voltages. For example, NAND flash memories use a polysilicon layer serving as a floating gate electrode layer, as a resistor.\nAs design rule dimensions decrease, the effect of, for example, a mutual interference resulting from the crosstalk between adjacent memory cell transistors becomes serious. Accordingly, to reduce the effect of such a mutual interference or the like, the floating gate electrode layer tends to have an increasingly reduced film thickness.\nThe reduced film thickness of the floating gate electrode layer significantly varies a resistance. The resistors elements may thus fail to provide appropriate functions, affecting circuit operations.\nFor example, a nonvolatile semiconductor memory has already been disclosed in which a gate oxide film material and a first gate electrode material (floating gate electrode layer) are sequentially deposited on a semiconductor substrate in this order, then an inter-gate insulating film material and a second gate electrode material (control gate electrode layer) are sequentially deposited on the first gate electrode in this order, and the second gate electrode material (control gate electrode layer) is formed into a resistor in a peripheral circuit portion (see, for example, Jpn. Pat. Appln. KOKAI Publication No. 2002-110825).\nWhen the control gate electrode layer is used for the resistor, the resistor disadvantageously has an increased area if its resistance remains unchanged. This is because the control gate electrode layer comprises a stack structure of silicide with metal such as tungsten or cobalt and polysilicon and thus has a lower resistivity than the floating gate electrode layer."}
{"text": "Roadside construction projects must oftentimes be carried out without stopping roadway traffic. Accordingly, measures are required to maintain safety in the work environment. Temporary warning signs are becoming increasingly popular since the geographic configuration of a work site frequently changes before the work is finished. Also, with the advent of more powerful, complex equipment, roadway repairs and roadside work projects can be completed within a relatively short time, such that erection of permanent warning signs would not be practical.\nOne class of temporary warning signs includes sign panels made of flexible material, which are maintained in a taut, planar position by backing members extending across the sign panel. These backing members frequently comprise ribs of fiberglass or aluminum material. It has become popular to arrange the sign panel so that it can be quickly and easily detached from the supporting ribs and rolled up for compact storage. Thus, the need arises for releasably securing the rib members to the flexible sign panel.\nInitially, pockets for receiving the ends of supporting ribs were formed by joining relatively small pieces of fabric to the corners of the flexible sign panels, so as to form a rib-receiving pocket between the fabric and the sign panel. Roadway warning signs are employed out of doors and are thereby subjected to wind gusts and traffic induced wind bursts which expose the flexible sign panel to substantial amounts of strain, particularly where rib supports are attached to the sign panel. Accordingly, with static and vibratory wind loadings applied to the sign panel, the ends of the supporting ribs were found to wear through the pocket material, requiring substantial repairs to be made to the sign assembly. In time, steps were taken to provide more secure attachment of supporting ribs to portions, usually outside corners, of the sign panel. For example, fabric pockets have been employed in U.S. Pat. Nos. 3,899,843 and 3,526,200, while resilient straps have been proposed in U.S. Pat. No. 4,592,158. Rigid panel pockets have also been proposed in U.S. Pat. Nos. 4,426,800 and 4,490,934. Despite these advancements, improvements in rib supports are still being sought."}
{"text": "Advances in minimally invasive surgical technology could dramatically increase the number of surgeries performed in a minimally invasive manner. Minimally invasive medical techniques are aimed at reducing the amount of extraneous tissue that is damaged during diagnostic or surgical procedures, thereby reducing patient recovery time, discomfort, and deleterious side effects. The average length of a hospital stay for a standard surgery may also be shortened significantly using minimally invasive surgical techniques. Thus, an increased adoption of minimally invasive techniques could save millions of hospital days, and millions of dollars annually in hospital residency costs alone. Patient recovery times, patient discomfort, surgical side effects, and time away from work may also be reduced with minimally invasive surgery.\nThe most common form of minimally invasive surgery may be endoscopy. Probably the most common form of endoscopy is laparoscopy, which is minimally invasive inspection and surgery inside the abdominal cavity. In standard laparoscopic surgery, a patient's abdomen is insufflated with gas, and cannula sleeves are passed through small (approximately ½ inch) incisions to provide entry ports for laparoscopic surgical instruments. The laparoscopic surgical instruments generally include a laparoscope (for viewing the surgical field) and working tools. The working tools are similar to those used in conventional (open) surgery, except that the working end or end effector of each tool is separated from its handle by an extension tube. As used herein, the term “end effector” means the actual working part of the surgical instrument and can include clamps, graspers, scissors, staplers, and needle holders, for example. To perform surgical procedures, the surgeon passes these working tools or instruments through the cannula sleeves to an internal surgical site and manipulates them from outside the abdomen. The surgeon monitors the procedure by means of a monitor that displays an image of the surgical site taken from the laparoscope. Similar endoscopic techniques are employed in, e.g., arthroscopy, retroperitoneoscopy, pelviscopy, nephroscopy, cystoscopy, cistemoscopy, sinoscopy, hysteroscopy, urethroscopy and the like.\nThere are many disadvantages relating to current minimally invasive surgical (MIS) technology. For example, existing MIS instruments deny the surgeon the flexibility of tool placement found in open surgery. Most current laparoscopic tools have rigid shafts, so that it can be difficult to approach the worksite through the small incision. Additionally, the length and construction of many endoscopic instruments reduces the surgeon's ability to feel forces exerted by tissues and organs on the end effector of the associated tool. The lack of dexterity and sensitivity of endoscopic tools is a major impediment to the expansion of minimally invasive surgery.\nMinimally invasive telesurgical robotic systems are being developed to increase a surgeon's dexterity when working within an internal surgical site, as well as to allow a surgeon to operate on a patient from a remote location. In a telesurgery system, the surgeon is often provided with an image of the surgical site at a computer workstation. While viewing a three-dimensional image of the surgical site on a suitable viewer or display, the surgeon performs the surgical procedures on the patient by manipulating master input or control devices of the workstation. The master controls the motion of a servomechanically operated surgical instrument. During the surgical procedure, the telesurgical system can provide mechanical actuation and control of a variety of surgical instruments or tools having end effectors such as, e.g., tissue graspers, needle drivers, or the like, that perform various functions for the surgeon, e.g., holding or driving a needle, grasping a blood vessel, or dissecting tissue, or the like, in response to manipulation of the master control devices.\nSome surgical tools employ a roll-pitch-yaw mechanism for providing three degrees of rotational movement to an end effector around three perpendicular axes. At about 90° pitch, the yaw and roll rotational movements overlap, resulting in the loss of one degree of rotational movement."}
{"text": "The present invention relates to a method for rotating at least two-dimensional image records with non-isotropic topical resolution.\nKnown literature contains a number of methods for rotating and shifting image records. These methods can be roughly divided into two categories. In the first category, the rotations and translations of a record are computed exclusively in the image space. In the second, the rotations and translations are computed completely or partially in Fourier space with the aid of Fourier transformation.\nThe methods of the first category are disclosed in the following examples: Starting with an existing image record, a rotated image record is obtained by executing the following two steps: In the first step, the coordinate points of the rotated image record are computed by multiplying the coordinates of the existing image record by a rotation matrix. The value of each rotated coordinate pointxe2x80x94this can be, for example, a matter of the percent value on a gray scalexe2x80x94is calculated in an interpolation process. A large number, more or less, of values of neighboring points of the existing image record, or all values, should be integrated into the interpolation process. The interpolations can be based on a Sinc interpolation, for instance. The following example should demonstrate that said interpolation processes entail a significant computing outlay. Assume a two-dimensional image with 256 pixels in two dimensions. This means that 65536 pixels must be rotated. In turn, this means that in the worst case for each of the 65536 pixels a respective interpolation must be performed with the same number of pixels. Altogether, this means a quite significant computing outlay and thus time outlay. For shifts of image records in accordance with methods of the first category, the steps described above for a rotation are performed accordingly.\nThe methods of the second category are exemplified next. There, Fourier transformation is specifically used to reduce the computing outlay. The Fourier transform F(k) of a one-dimensional image space function f(x) is generally defined as       F    ⁡          (      k      )        =            ∫              -        ∞                    +        ∞              ⁢                            f          ⁡                      (            x            )                          ·                  ⅇ                      j            ·            2            ·            π            ·            k            ·            x                              ⁢                        ⅆ          x                .            \nFor the inverse-transformation,       f    ⁡          (      x      )        =            ∫              -        ∞                    +        ∞              ⁢                            F          ⁡                      (            k            )                          ·                  ⅇ                                    -              j                        ·            2            ·            π            ·            k            ·            x                              ⁢                        ⅆ          k                .            \nAs the central feature of the Fourier transformation, the displacement set is of primary importance in the methods of second category. The displacement set for a one-dimensional Fourier transformation is as follows:\nf(x-x0) in the image space corresponds to ejxc2x72xc2x7xcfx80xc2x7x0xc2x7kxc2x7F(k) in Fourier space.\nApplied to multidimensional applications, this means that a shift of an image record by a random vector is imaged in Fourier space in an additional phase of the Fourier transform. Unlike methods of the first category, interpolations are not required for the translation of an image record. The interpolation is completed on the basis of the attributes of Fourier transformation semi-automatically with the transformation into Fourier space, with the introduction of a desired additional phase and the corresponding back-transformation. Rotations of image records can likewise be performed with the aid of Fourier transformation and its displacement set without interpolation. To accomplish this, a rotation matrix describing the rotation is dismantled into a product of corresponding shear matrices. The effect of shear matrices is expressed in simple shifts of rows or columns of an image record, respectively. These shifts can be calculated easily in Fourier space with the aid of the Fourier transformation and its displacement set. A method of another kind is described in the essay by William F. Eddy, Mark Fitzgerald and Douglas C. Noll, xe2x80x9cImproved Image Registration by Using Fourier Interpolationxe2x80x9d, MRM 36, pp. 923-931, 1996.\nThe case of a non-isotropic topical resolution, which is the more frequent case in many applications, can be handled by the above method of the second category only with additional outlay. Non-isotropic records are not problematic for methods of the first category. For example, in rotations an elliptical path is merely executed instead of a circular path. The high computing outlay for the necessary interpolations remains unchanged. In methods of the second category, the non-isotropic record must be converted into an isotropic auxiliary record in a first step. Then, a method of the second category can be applied. Finally, the record must be converted back into a non-isotropic record. These conversion processes, which are known as resampling processes, represent additional computing outlay, regardless of the type and manner of execution of the resampling processes. In practice, the registering of non-isotropic records is often avoided for these reasons.\nIt is thus an object of the invention to create a method which reduces the computing outlay in rotations and translations of non-isotropic image records.\nThis object is inventively achieved in accordance with the present invention in a method for rotating at least two-dimensional image records with non-isotropic topical resolution, said method comprising the steps of: describing a rotation of an image record using a rotation matrix; representing said rotation matrix as a product of at least two shear matrices, each shear matrix having at least one element that is dependent on the angle of said rotation and remaining elements which are exclusively zeroes and ones; multiplying a matrix element, which is dependent on said angle of said rotation, of at least one of said shear matrices by a factor; and performing said rotation of said image record in Fourier space without interpolations and without forming an isotropic auxiliary record, upon exploitation of a displacement set of the Fourier transformation by implementing said shear matrices as displacements of line elements of said image record.\nIn an embodiment, the present method provides the following steps:\ndescribing a rotation of an image record with a rotation matrix;\nrepresenting the rotation matrix as a product of at least two shear matrices, each of which comprises exactly one element that is dependent on the angle of rotation, and whose remaining elements are exclusively zeroes and ones;\nmultiplying the matrix element, which depends on the angle of rotation, of at least one shear matrix by a factor; and\nperforming the rotation of the image record in Fourier space, without interpolations and without forming an isotropic auxiliary record, by exploiting the displacement set of the Fourier transformation by implementing the shear matrices as displacements of line elements of the image record.\nBy simple one-time multiplications, by factors, of matrix elements of shear matrices that describe rotation intensive transformations of non-isotropic image records into isotropic auxiliary image records, corresponding back-transformations are made superfluous for the purpose of rotating non-isotropic image records in Fourier space. The application of the invention in a computer system reduces the necessary computing time.\nIn an advantageous embodiment, the factor depends exclusively on at least one ratio of a topical resolution in a first dimension to the topical resolution in a second dimension, which ratio describes a non-isotropic topical resolution. With the aid of a factor that is selected in this way, the error between the image record that is rotated as desired and the image record that is rotated by applying the method is minimized.\nIn a particularly advantageous embodiment, said error equals zero for two- or three-dimensional image records.\nAn advantageous embodiment relates to image records that are generated by a computer or MR tomography device. Here, the present invention can be advantageously employed particularly in the field of functional MR tomography, where the image records that are registered in a time series are rotated and displaced continuously for the purpose of motion correction."}
{"text": "1. Field of the Invention\nThe present invention relates to a memory driving method and a semiconductor storage device, particularly to an FBC (Floating Body Cell) memory in which information is stored by accumulating majority carriers in a floating body of a field effect transistor.\n2. Related Art\nRecently, there is an FBC memory device as a semiconductor storage device expected to be an alternative to a 1T (Transistor)-1C (Capacitor) type DRAM. The FBC memory device is formed by FET (Field Effect Transistor) including a floating body (hereinafter referred to as body) on an SOI (Silicon On Insulator) substrate. In the FBC memory device, data “1” or data “0” is stored according to the number of majority carriers accumulated in the body. For example, in FBC formed by an N-type FET, it is conventionally defined that the data “1” is the state in which many holes are accumulated in the body while the data “0” is the state in which few holes are accumulated in the body. A memory cell in which the data “0” is stored is called a “0” cell and a memory cell in which the data “1” is stored is called a “1” cell.\nFBC is excellent in miniaturization compared with a conventional DRAM. However, because the body of FBC has an electrostatic capacity smaller than that of a capacitor of the conventional DRAM, FBC has a data retention time shorter than that of DRAM although a leakage current from the body of FBC is smaller than that from the capacitor of DRAM. Accordingly, it is necessary to frequently perform a refresh operation. As a result, disadvantageously a ratio of a time when usual read and write are prohibited (refresh busy rate) is increased and a current necessary to retain the data is increased compared with the conventional DRAM. In particular, large power consumption becomes a serious problem in portable devices.\nAdditionally, in the FBC memory, it is necessary that a size of a current driver is enlarged because a current is passed through the memory cell to write the data. Therefore, a chip size is not so decreased although the memory cell has the small size. That is, a proportion of the memory cell to the chip (cell efficiency) is small.\nThere has been proposed block refresh in order to deal with the problems (P. Fazan, S. Okhonin and M. Nagoga, “A new block refresh concept for SOI floating body memories” IEEE Int. SOI Conference, pp. 15-16, September, 2003). The block refresh is a method, in which holes are supplied only to the “1” cell by impact ionization and the holes are drawn from both the “0” cell and “1” cell by utilizing a charge pumping phenomenon. The charge pumping phenomenon is one in which the holes are drawn from the body as a result of recombination of the holes in the body and electrons trapped by a surface state existing in an interface between a silicon substrate and a gate dielectric film. Accordingly, density of the surface state becomes important. Usually the surface state has the density of about 1010 cm−2. For example, an average of about one surface state exists for a channel having an area of 0.1 μm×0.1 μm. That is, the memory cell having no surface state exists with a considerably high probability. Thus, the block refresh is not the effective method to the memory cell having no surface state, and the block refresh is impractical."}
{"text": "1. the Field of the Invention\nThis invention relates to sound generation and, more particularly, to novel systems and methods for providing aural and tactile signals to a human being.\n2. The Background Art\nAutism is a condition characterized by persons being overwhelmed by sensory perceptions, which become distorted (disoriented). Typically, the life of an autistic individual is characterized by overwhelming chaos due to sensitivity to external stimuli, including sight, sound, touch, taste, and smell. Various systems have been developed for treating autism.\nPrevious work has been done by the inventor in the area of orientation counseling and procedures. It has been provided in the past to place a headset or earbuds on a user in order to generate a sound that is detected by the ears, processed by the brain, and treated as a focal point or an orientation point for an individual. Once focused, perceptions become accurate or undistorted (oriented). However, many persons, particularly young children or severely affected subjects, may not be able to wear a headset in order to listen to the sounds that will help to focus attention, direct the “mind's eye,” toward the sound perceived to be at a particular position in space.\nFor example, the essence of stereo is the stereophonic projection of sound from different locations. Stereophonic headphones, for example, permit different tracks of music or other sound to be passed to the ears, thus giving an impression of position and distance. In reality, stereo sound is propagated by a speaker relatively nearer each ear. However, by changing the volume of different tracks, particular sounds may be played back to a listener in such a way as to appear to be originating at different locations, to move with time, or the like. However, experts on the subject have informed the inventor that it is impossible to exactly locate in space a perceived location for the sound. More particularly, the inventor was informed that it was impossible to maintain such a location with respect to an individual.\nThe David Autism Approach is documented in various books and websites available internationally. Unfortunately, many autistic individuals are unable to wear headphones, or unwilling to do so. Thus, an ability to receive a sound signal, which has been generated with certain artifacts (such as accommodating the geometry and attenuation of the hearer's head) and balanced to seem like it is coming from a particular point in space fixed with respect to the person, is not heretofore possible. However, it would be an advance in the art to provide a system, that generates sound, in a way that such a system can be created and then worn in such a location and manner that it is not easily removed by a subject, and does not interfere with movement of the head, and normal daily activities."}
{"text": "Characteristic of cardiothoracic surgery is the post-operative patient who is sent to the Intensive Care Unit (ICU) intubated due to respiratory requirements. Approximately half of these patients are extubated within their first twenty-four post-operative hours. In most cases these patients are extubated within the first three days. There are some, however, who remain intubated for a significant length of time. When a surgeon identifies a patient who requires intubation longer than seven days, the surgeon will usually decide to perform a tracheotomy on that patient. The breathing support tube enters the trachea rather than entering the mouth for the trached patient. Communication for a intubated or trached patient is minimal due to the inability to speak resulting in the patient, hospital staff and loved ones resorting to the reading of lips, nodding of heads and squeezing of hands to communicate.\nWithout effective communication, the intubated or trached patient may not receive the standard of care he or she would otherwise receive had he or she been able to effectively communicate. The lack of communication also creates unnecessary levels of anxiety which the patient must endure. Nurses and hospital staff ask many questions from the patient pertaining to their prognosis and progress which may never get fully or even adequately answered. A doctor or nurse is not able to treat a symptom which they know little or nothing about. In addition, other problems arise due to the insufficient communication from the patient. Localized areas of pain are often mis-diagnosed, resulting in over-medication generally or the medication of an area which is not the source of pain. Proper and essential treatment given in an adequate and timely manner will help resolve or prevent many post-operative complications and decrease the patient's length of stay in the hospital. This begins with providing the patient a clear and precise means of communication.\nAnother problem exists in that currently patients are subjected to pushing a button or call light, which turns on a light in the hallway at the doorway to their room. Nurses have no way of identifying whether the patient's need is urgent or non-urgent. Additionally, the nurse is unable to prepare him/herself for the need appropriately before entering the room. Instead, the nurse must go to the patient's room, be informed of the problem or need and then leave the patient's room and retrieve whatever resources are necessary for the nurse to fulfill the patient's need or request. This is extremely time-consuming, wastes precious hospital resources, and can delay meeting patient's needs. This problem can be detrimental to the patient when the need is of an urgent matter. Unless the patient can scream loud enough to be heard from wherever help may be, the patient is subjected to wait until someone responds to a common light at the patient's doorway.\nMoreover, current systems lack a patient-centric device for the bedside interface (e.g., pillow speaker); and, only a few nurse call systems provide an opportunity for patients to convey specific messages directly to their assigned providers.\nFurthermore, eighty percent of hospitals care for patients with limited English proficiency (LEP) on a regular basis, and despite advancements in the profession of healthcare interpreting and translation services (aka, Language Access Services), patients with language barriers are often left without an effective means to communicate with their providers. While best practice, clinical ethics and legal and regulatory guidelines recommend the use of professional interpreters for all healthcare encounters, logistics and resource capacity make this prohibitive. Reasons cited by hospital staff for not using professional interpreters include resources available to bridge the language barrier are not user friendly; resources are not easily accessible; staff are unaware of the resource and have not been trained. Despite these disadvantages, nurse call systems do not provide a means for LEP patients to generate a nurse call request in the patient's preferred language.\nAccordingly, it has been estimated that inefficient communication costs U.S. hospitals more than $12 billion annually or $4 million for each 500 bed hospital. In summary, nurse call systems have been the primary means for hospitalized patients to initiate an encounter from the bedside. However, these nurse call patient requests range in urgency, are not differentiated based on skill-set required to fulfill the patient request, and are not equitable for LEP patients. Further, these shortcomings prohibit effective communication with LEP patients and can contribute to poorer outcomes compared to their English-speaking counterparts."}
{"text": "A portable exerciser is needed to stimulate blood circulation while riding in a car or an airplane or while confined to bed, but no such exercise device is presently known which is as simple in operation and as effective as the present invention."}
{"text": "High volume, high resolution printing is an objective that has been sought by the manufacturers of wide format printers for some time. Wide format printers have been available to the public for many years. Examples of popular wide format printers are the Hewlett Packard (HP) 1000/5000, the HP 3000/3500, the Epson 7000/10 000 and many others.\nThese printers all have a traversing printhead that traverses a print medium while depositing ink on the medium. Applicant believes that these printers suffer from inherent disadvantages, particularly when attempts are made to utilize the design of such printers in order to achieve faster printing speeds at high resolutions.\nCentral to the problem of achieving high printing speeds is the ability to achieve a printhead that is capable of generating the necessary number of ink dots at a suitable rate. Further, in order to achieve accurate printing, it is desirable that a row or band of the image be created in as little print cycles as possible, and preferably in a single print cycle. It follows that it is undesirable for a traversing printhead to be used in an attempt to achieve high print speeds and that a single printhead incorporating a suitable number of inkjet nozzles is required.\nThermal printheads also referred to as bubble jet printheads and piezoelectric printheads have been available for some time. These suffer from excessive heat build up and energy consumption and have therefore been found by the applicant to not be suitable for use in a pagewidth configuration. A number of disadvantages associated with such printheads are set out in U.S. application Ser. No. 09/526,504.\nThe applicant has developed a printhead chip that is capable of producing images having a resolution as high as 1600 dpi. These chips are manufactured using integrated circuit fabrication techniques. Details of the chips are provided in the above referenced applications and patents. Applicant believes that these printhead chips are extremely suitable for use in wide format printers. The reason for this is that such chips operate at extremely high speeds due to the large number of nozzle arrangements required in a single chip and due to the fact that such chips can be driven at an extremely high cyclical rate.\nThe Applicant has been faced with a number of difficulties in order to achieve the effective use of such printhead chips in wide format printers. One particular difficulty identified by the Applicant is the effective control of a number of such printhead chips to achieve accurate printing. This control must incorporate the use of effective image processing tools that are capable of processing stored images at a rate that corresponds with the physical rate of printing achievable by a number of the above printhead chips.\nAnother difficulty that faces the manufacturers of wide format printers are the problems associated with heat build up. This can often result in the necessity for expensive heat extraction devices that add to the complexity of the printer."}
{"text": "1. Field of the Invention\nThe present invention relates to an address translation control system, more particularly to an address translation control system which, in addition to an address translation buffer performing a high speed translation from a logical address to a physical address, comprises a memory array copied from a portion of the contents of the address translation buffer, and in which a partial purge in the address translation buffer is performed by searching the memory array.\n2. Description of the Related Art\nIn a data processing system using a virtual storage method, the translation from a logical address to a physical address at a high speed is performed by an address translation control system wherein address translated pairs of logical addresses and physical addresses are registered in a TLB (Translation Lookaside Buffer), i.e., an address translation buffer, and a high speed address translation is performed by accessing the address translation buffer.\nThe address translation buffer comprises a primary block and an alternate block, wherein each block has a plurality of entries. Each entry comprises a valid flag showing the validity of the entry, a multiple virtual storage identification field identifying a multiple virtual storage space, a logical address field, a physical address field, and a storage protection key field.\nIn a conventional address translation control system having such an address translation buffer, for an ordinary command fetch or operand access, the above-mentioned address translation buffer is accessed by applying a portion of the logical address to be translated, i.e., access address from an effective address register. By this accessing, the entire contents of the accessed entry stored in the primary block and the alternate block in the address translation buffer are read. A multiple virtual storage identification and the partial logical address which are read from the address translation buffer are compared with a multiple virtual storage identification from a multiple virtual storage identification register and another portion of the logical address from the effective address register, respectively. If the identifications of the comparing contents are obtained, the physical address in the entry is sent to a real address register as a physical address corresponding to a logical address in the effective address register.\nIn a virtual storage system, since the logical address covers a wider range than the utilizable physical address, a physical address, in turn, is allocated to another logical address when the physical address becomes non-utilizable. At that time, if a translation table in a main storage unit is renewed, since pairs of logical addresses and physical addresses translated before the renewal may remain in the address translation buffer, the pairs must be searched for and then purged. For this reason, in a conventional address translation control system, a partial purge method wherein each entry in the address translation buffer is searched for a physical address, is frequently used. According to this method, when a physical address to be purged is applied to a purge register, the access address in the effective address register is sequentially increased, an entry in the address translation buffer is sequentially read, the physical address in the purge register is compared with the physical address in the address translation buffer, and if the addresses are identical, a valid flag of the entry is made invalid.\nIn such a conventional address translation control system, when the partial purge is carried out, since all the contents of the address translation buffer must be read out, there is a defect in that the address translation buffer cannot be used at an ordinary command fetch or operand access.\nThe above defect in the conventional address translation control system causes a decrease in the operating efficiency of the address translation buffer, and in particular a considerable reduction in the performance when using the pipeline method."}
{"text": "1. Field of the Invention\nThe present invention is in the field of vehicular power plants and more specifically relates to a type of power plant in which an internal combustion engine operates a generator that charges a battery pack, while simultaneously, a steam engine using waste heat from the internal combustion engine operates a second generator that also charges the battery pack. There is no mechanical connection between either engine and the drive wheels, the drive wheels being driven by electric motors powered by the battery pack.\n2. The Prior Art\nIn U.S. Pat. No. 4,405,029 issued Sept. 20 1983, Hunt shows a vehicle that has an electric motor as well as an internal combustion engine that may selectively be connected to drive the wheels of the vehicle, either independently or together. At speeds above 20 mph, the internal combustion engine drives the wheels, and heat from the engine is used to generate electric power that is stored in the battery. Below 20 mph, the wheels are driven by an electric motor powered by the battery, although the internal combustion engine may be used to assist. This differs from the present invention in that, in the present invention the battery-powered electric motor always drives the wheels and the gasoline engine never directly drives the wheels but only operates a D.C. generator.\nIn U.S. Pat. No. 4,470,476 issued Sept. 11, 1984, which is a divisional of the aforementioned U.S. Pat. No. 4,405,029, Hunt shows another embodiment of his invention, in which radiant energy of the waste heat and exhaust gases from the internal combustion engine are converted directly to electrical energy by means of photovoltaic cells which are used to charge the battery. In contrast, the present invention does not use photovoltaic cells.\nIn U.S. Pat. No. 4,109,743, Brusaglino et al. describe a vehicle propulsion system in which a turbine drives two counter rotating flywheels that in turn drive two D.C. generators. These D.C. generators power electric motors that turn the drive wheels. In this system the energy is stored in the flywheels rather than in a battery pack, and no batteries are used in the system.\nIn U.S. Pat. No. 4,394,582, Kreissl et al. show a method and apparatus for using the waste heat energy from an internal combustion engine. In their system, a turbine is driven by exhaust gases from the internal combustion engine, and the turbine drives a generator. The generator connected through an electrical converter to an electric motor which is drivingly coupled to the internal combustion engine to help it to overcome its load. The turbine needs to turn very fast (120,000 rpm) and so it is necessary to use a \"synchronous\" type of rotor. Thus, Kreissl et al. use waste heat as a way of mechanically assisting the internal combustion engine that is connected to the drive wheels of the car.\nIn U.S. Pat. No. 4,570,077 issued Feb. 11, 1986 to Lambley, there is shown a waste heat recovery system in which the waste heat from the main engine heats a boiler that drives a turbine that is used to drive an alternator. When there is not enough waste heat, a hydraulic drive powered by the main engine is clutched to the alternator rotor shaft to turn it.\nIn U.S. Pat. No. 4,433,548 issued Feb. 28, 1984 to Hallstrom, Jr. there is described an engine that is a combination internal combustion engine and steam engine, and in which the hot exhaust gases resulting from the power stroke of the internal combustion engine are used to generate steam that is used to provide a second power stroke.\nIn the latter three patents described above, the wheels of the automobile are driven by an internal combustion engine, while in contrast, in the present invention the wheels are always driven by a D.C. drive motor powered by a battery. The battery, in turn, is charged by a first D.C. generator powered by an internal combustion engine and by a second D.C. generator powered by a steam engine that is heated by waste heat from the internal combustion engine."}
{"text": "A Fabry-Perot interferometer is an instrument commonly used in high resolution optical spectroscopy, and as a tool for the analysis of laser radiation. Etalons are widely used in telecommunications, lasers, and spectroscopy for controlling and measuring the wavelength of light. For example, a laser beam typically comprises several discrete optical frequencies related to different modes of oscillation of the laser resonator. The Fabry-Perot interferometer can be used to determine the number of modes oscillating.\nThe Fabry-Perot interferometer comprises two parallel mirrors and a Fabry-Perot etalon which has certain optical resonator properties. One of the mirrors is fabricated at different distances (two or more levels) from the other mirror. Conventional spectrometers comprise a large, expensive, and complex system of lenses and mirrors. More recently, attempts have been made to manufacture a Fabry-Perot interferometer within a monolithically integrated circuit. See for example, “Single-chip CMOS optical microspectrometer”, J. H. Correia et al., Sensors and Actuators, 82 (2000) 1.91-197.\nThese etalon fabrication of fixed cavity filters in monolithically integrated circuits include multiple timed etches for providing varying discrete levels of one mirror. However, establishing the distance of the mirror levels by a timed etch is inaccurate for dimensions involving wavelengths on the order of 400 to 1000 nanometers.\nAccordingly, it is desirable to provide a method of fabricating a Fabry-Perot filter in an integrated circuit whose mirror levels may be accurately determined. Furthermore, other desirable features and characteristics of the present invention will become apparent from the subsequent detailed description of the invention and the appended claims, taken in conjunction with the accompanying drawings and this background of the invention."}
{"text": "The invention relates to a cleaning sieve for the harvested crop of agriculatural machine and in particular harvester threshers.\nIn present-day modern harvester threshers, two sieves one disposed above the other are generally employed wherein the upper one is shaped as a lamellar sieve which separates the coarse chaff and straw particles and the lower one sieves the finer remaining particles. The two sieves are passed by the air current from a blower from below.\nThe air current lifts the lighter chaff and straw particles so that they are transported by the air current above the sieves and the vibratory movement of the sieves to the rear of the machine and out of it whereas the specifically heavier grains fall through the sieves and are collected in a manner known per se in a grain tank or are filled into sacks.\nIn modern harvester threshers, the sieves are longer than one meter. Chaff and short straw have to be transported over the whole length of the sieves. The air current is not in this case able to keep the chaff particles coming up at the beginning of the sieve floating over the whole length of the sieve and to transport them toward the rear. Otherwise the air current should be so strong that it could also keep grains afloat and blow them out of the machine together with the chaff.\nIt has therefore been tried again and again by using mechanical means to keep the chaff and straw particles as far as possible above the upper sieve and contrary thereto conduct the heavier grains from the very beginning to the upper sieve in order to separate grains and chaff at least partly from each other before they fall onto the upper lamellar sieve passed by the air current.\nIn German Patent Specification No. 908,080, a structure of this sort has been described wherein on the gliding bottom from which the grain-chaff mixture is conducted onto the upper sieve, a rake-like extension is provided which keeps the chaff and straw particles over a certain height above the upper sieve while the grains drop through the rake onto the upper sieve.\nA nearly identical apparatus has been described in GDR Patent Specification No. 227,604. In order to loosen the chaff-grain mixture on the upper lamellar sieve, most various sieve structures have been suggested in the past. German Utility Model No. 7,632,151, for instance describes a ramp sieve which is similar to a lamellar sieve in function but has the disadvantage as compared thereto that it comprises square hole openings instead of slit openings as in case of the lamellar sieve.\nSuch sieves were not successful since an adjustment of the sieve openings, which was obviously necessary with a view to the various kinds of fruit and humidity grades, is not possible in most cases and, secondly, longer straw particles passing through the holes and slipping into the lower sieve will stick fast and after a while will block the whole sieve.\nIn order to loosen the grain-straw mixture on a lamellar sieve, it has been suggested in German Patent Specification No. 1,238,707 to provide individual lamellae with finger-shaped extensions similar to a rake whereby these lamellae are vibratory relative to the other lamellae, whereby these finger-shaped extensions perform an up and down movement relative to the remaining lamellae and so loosen up the chaff-grain mixture.\nThis structure was not successful since the vibrating lamellae, on one hand, perform an on-off movement relative to the remaining sieve lamellae, i.e. do not permit a fixed preadjusted value corresponding to the kind of fruit to be harvested, and the wear of the pivotably moving lamella axes, considering a frequency of a couple of hundred movements per minute, leads to a breakdown of the sieve after a short time, on the other.\nIn German Offenlegungsschrift No. 2,751,842, a rake-like apparatus disposed above the sieve and provided particularly for the corn harvest has been described which however is completely unsuited for loosening a grain-straw mixture considering that the straw and chaff particles dropped under the rake prongs are kept down by these very prongs and hence cannot be pressed upward by the air current. A sieve provided with such an apparatus will be clogged within a short time.\nU.S. Pat. No. 4,548,213 describes a sieve also particularly developed for the corn harvest which is inserted without a lower sieve and on which, on one hand, ribs are provided which laterally support the material gliding over the sieve when in inclination, and ondulated transversal ribs, on the other hand, prevent that the corn cobs glide too fast over the sieve and are lost. Such a sieve is not suited for cereals since the effect of a lamellar sieve for controlling the air current is lacking.\nIn German Patent Specification No. 3,720,696, a rake-like cover device provided above the oversweep area of a sieve box has been described which is to prevent that the grain-chaff mixture collecting on one side, when in an inclined position, is guided in total to the oversweep device of the harvester thresher. This cover is said to effect a last separation of the grains from the mixture. The performance of the sieve above the working zone as such, i.e. before the oversweep area, is not in this case influenced nor is it increased. A similar apparatus has been described in German Offenlegungsschrift No. 4,106,814 wherein, above the oversweep area, enlarged lamellar intermediate spaces are provided which when in an inclined position will guide the collecting grain-chaff mixture to the oversweep. This apparatus, is not in any way suited either to increase the performance of the sieve above the working zone as such.\nIn both the apparatus according to German Patent Specification No. 3,720,696 and German Offenlegungsschrift No. 4,106,814 referred to above the very opposite effect can be experienced. The oversweep amount obtained in excess because of the larger passages in the oversweep area starts a second trip over the sieve device and increases the crop amount on the sieves to be processed. This can lead to a situation where after a second trip, a third and even multiple trip develops which will overload the sieve to an extent that the machine has to run idly before mowing can be resumed.\nThe sieve apparatus most effective for the time being has been described in German Offenlegungsschrift 3,704,348. This apparatus comprises a plurality of separated upper sieves provided on different levels. By this arrangement, drop steps are provided between which the air current blows the chaff and straw particles towards the rear portion, in the travelling direction, of the following sieve provided on a lower level. With a view to this arrangement, the dreaded \"carpet\" which blocks the air current through the sieve cannot develop. The disadvantage of this sieve arrangement however is to be seen in that a vertically much higher space requirement of the sieve box is necessary. The harvester thresher has to be constructed from the very beginning to hold such a sieve arrangement. For space reasons, it will almost never be possible to add this arrangement to machines already available."}
{"text": "This invention relates to a lit-up marking device of the type which has an angular section, used for positioning on the front vertex of the step, which defines two substantially perpendicular branches and a longitudinal ducting for fitting means of lighting and a translucent or transparent protector.\nAs background of the invention, Utility Models 293.909, 9802069 and 9803048, held by the applicant for this invention, may be mentioned.\nUtility Model 293.909 relates to a protector for the front edge of the steps which is made up of an angular section whose front, upper zone has a continuous chamber for the fitting of a string of incandescent lights and a transparent or translucent plate which acts as protector for the lights which form the step-marking means.\nUtility Model 9802069 describes a marking device for floors which, having an arrangement similar to that of the Model described above, has the particular feature that the ends of the wings of the section are slightly raised to form carpet retainers. The upper wing of the section also presents a longitudinal seat for fitting a strip of anti-slip material, thereby preventing users slipping when they walk on the metal section. Said Model further includes some wedge-shaped sections which act as a trim for the lit-up marking section where no carpet has been laid on either of the sides of same.\nUtility Model 9803048 relates to improvements in lit-up marking device for grandstands, steps and the like; in accordance with said improvements, the section for fitting onto the front side of the step, grandstands and the like has a generally L-shaped layout whose front wing is formed by a vertical prolongation which runs from the longitudinal channel in which the lighting means are housed.\nThis wing has an end portion with which it makes contact with the front of the step, so that its middle zone is separated from the latter to form a cable ducting. The aforesaid front wing has on its front surface longitudinal guides for fitting an advertising plate or support. This advertising plate or support is lit from the back zone, for which purpose some orifices have previously to be made in the section for fitting of the corresponding lighting means in a position facing them.\nThis Model also makes provision for fitting a row-numbering plate on one of the side ends of the upper seat, by using the remaining length of the upper portion of the seat for fitting the anti-slip strip of material mentioned above.\nThe row-numbering plates are lit up from the lower zone, for which purpose some holes have to be made previously in the upper wing of the section for fitting of the corresponding lighting means in a position facing them.\nSaid improvements, by providing novel characteristics, such as the possibility of including row-numbering plates or an advertising support, involve some problems for the manufacturer and the fitter, such as the need to keep a larger stock since, depending on the end customer, a section must be fitted which permits the attachment of advertising or no attachment thereof, or the need to drill in the section some holes for fitting the means of lighting of the row-numbering plate and the advertising support.\nPatent GB 2 115 451 describes: xe2x80x9ca unit for illuminating a step comprising a combination of a stair tread moulding and a second moulding, the second moulding being formed with a recess accommodating a lighting unit, and, in use, being a clamped adjacent a stair riser by a downward flange on the stair tread mouldingxe2x80x9d.\nThis illuminating unit does not include means for attaching advertising nor the possibility of adding same.\nAmong the background art of the invention we may also stress the existence of U.S. Pat. No. 5,340,627, relating to a step-illumination device comprising: xe2x80x9ca light fixture housing extrusion having first and second channels formed therein, and a vertical riser plate intersecting with a horizontal step plate at a right angle, the horizontal step plate including a top surface having an integrally-formed lip forming a first carpet insert slot; and a means for removably securing a light fixture in at a means one of the first and second channelsxe2x80x9d.\nPatent GB 2 115 451 describes: xe2x80x9ca unit for illuminating a step comprising a combination of a stair tread moulding and a second moulding, the second moulding being formed with a recess accommodating a lighting unit, and, in use, being a clamped adjacent a stair riser by a downward flange on the stair tread moulding.xe2x80x9d\nThis lit-up marking device for steps, which is of the type that includes a section for positioning same on the front edge of the step or grandstand in question, which section presents two substantially perpendicular wings, a recess in its exterior vertex for fitting lighting means and a protector element, and an upper seat for fitting a strip of non-slip material and a row-numberer has the special feature of including an auxiliary section provided with guides for the fitting of an advertising support which is provided backing onto the riser or vertical surface of the step to be marked, while both sections present complementary means for attachment and detachment of said sections under pressure.\nThis special feature permits, according to customer needs, fitting of solely the upper section with the lit-up marking means for the step or the grandstand, or the fitting, together with it, of the auxiliary section which includes the guides for fitting the publicity support.\nThe complementary means for pressure-fitting of the two sections, principal and auxiliary, are defined respectively on the rear side of the former and on the front side of the latter, so that in the fitting position they overlap partially, in such a way that at a glance they appear to form a single section.\nThe lower end of the front wing of the principal section has a longitudinal channel with guides for fitting of the lighting means for the advertising support and recesses for hooking on a transparent or translucent protector, so that it is not necessary to machine any of the sections in order to fit said lighting means.\nIn accordance with the invention, the guides for fitting the lighting means are spaced at a short distance from the bottom of the aforesaid channel, leaving defined between them a hollow section for running the power-supply cables.\nThe seat defined in the principal section for fitting of the non-slip strip and the row-numberer has on its ends some interior mortised inlets, one of which is deeper so as to allow the power-supply cables of the means of lighting of the row-numberer to be run inside it, without any need to machine the section.\nIn accordance with the invention, the row-numberers have on the side facing the deeper mortised inlet an orifice for housing the corresponding lighting means.\nThis invention also includes other constructional features which affect: the translucent protector, in order to make the lit-up marking easier on the eye; the means of lighting used to achieve uniform-lighting across the entire width of the section; the side arrangement of LEDs for lighting up the optical fibre, the row-numberers and the advertising supports, together with incorporation of some ducts for the running through of cables for the various steps or the grandstand tiers and for concealed fitting of the LEDs or microbulbs used for side-lighting of the various elements.\nAs mentioned immediately above, the lighting means can be composed of a section of optical fibre or a slightly shorter length equal to the width of the step, for fitting the longitudinal recess of the section and some LEDs or microbulbs facing at least one of the ends of the section of optical fibre.\nOptionally, the LEDs or microbulbs which light up the optical fibre can be housed in one of the ends of the translucent protector or in orifices made in an angle piece, of opaque material, for fitting onto the side ends of the front edge of the step or grandstand, and forming part of ducting for the power-supply cables.\nWhere the LEDs or the microbulbs are housed in the end of the translucent protector, that protector will have on the corresponding end an opaque section which prevents it being seen by the user.\nThe angle piece will be used for those cases in which the steps do not have individual power sockets, with such sockets fitted only on the end steps.\nIn both cases, therefore, the advantage is that the user is not exposed to an intense point of light in the zone in which the LED is located, since the latter is covered either by the angle piece, which is opaque, or by the opaque section of the protector.\nThe angle piece has side appendages for lateral attachment thereof onto the section placed on the front angle of the step and appendages on its rear and lower ends for attachment to cable-carrying terminal strips inside a longitudinal duct which, together with the angle pieces, form ducting for running the power-supply cables through the successive steps.\nProvision has been made for the angle piece to have additional orifices for optional fitting of LEDs or microbulbs that side-light the row-numberers and the advertising supports in those cases in which the steps do not have individual power sockets and recourse has to be had to fitting the ducting mentioned above.\nSaid angle piece has a lower channel for passing cables from the terminal strips to the side orifices in which the LEDs or microbulbs are fitted.\nIn order to improve the marking of the steps, and in addition to the use of lengths of optical fibre lit up by the concealed LEDs, provision has been made for the translucent protector fitted on the longitudinal recess to have an upper portion coplanar with the upper surface of the section positioned on the front angle of the step, which makes it easier for users to see, especially when they are descending the stairway.\nIn accordance with the invention, the translucent protector has to the exterior a convex-curved surface which projects from the longitudinal recess of the section and an upper portion coplanar with the upper surface of the section which lights up a front portion of the horizontal surface of the step, thus making it easier to see for people descending the stairway.\nThe exterior surface of this translucent protector can have a longitudinal toothing that acts as a light diffuser and prevents users slipping if they tread on it.\nOptionally, this protector may have an open, U-shaped arrangement or a tubular arrangement, in this last case being prolonged by the upper zone of the section being extended into a horizontal lug which ends in a U-shaped piece for fixing into one of the upper mortised inlets of the section, and by the front zone of the section into a convex-curved lug whose middle zone projects into a convex-curved lug whose middle zone projects beyond the vertical defined by the front of the section."}
{"text": "1. Field of the Invention\nThe disclosed technology generally relates to single-carrier communication systems wherein adaptive feedback equalization is applied.\n2. Description of the Related Technology\nIn outdoor wireless applications the multipath channel exhibits very long impulse responses, which can lead to delay spreads of several tens of microseconds. This can be of particular importance in e.g. cellular or broadcast applications. When transmitting at high rate over such channels, the introduced intersymbol interference (ISI) can span hundreds of symbols and, hence, severely distorts the received signals. To cope with such ISI, very long decision feedback equalizers (DFE) are required at the receiver to properly recover the transmitted signal. The design of such equalizer can be very complex and occupy a large chip area since certain performance requirements have to be met.\nDecision feedback equalizers (DFE), consisting of a feedforward (FF) and a feedback (FB) filter, are well known in the art. They are preferred to linear equalizers because of their effectiveness at reducing the ISI. This stems mainly from their ability to cancel very efficiently the post-cursor portion of the ISI. This is done by optimally computing the feedforward and feedback coefficients, provided that the channel impulse response is known at the receiver. When this is not the case and the transmission extends over a large number of symbols (as e.g. in broadcasting systems), an adaptive DFE provides the means to compute adaptively the filter coefficients without any prior knowledge about the channel. According to the system definition, this adaptation can be carried out in a data-aided (trained) mode or non data-aided (blind) mode.\nAlthough trained or blind DFEs may provide satisfactory performance, even in high delay spread channels, their implementation complexity grows linearly with the channel length expressed in number of symbol instants. When this length reaches several hundreds of symbol instants, the number of operations can easily exceed a thousand complex multiplications per symbol. For a symbol rate of 10 MHz, which is typical of broadband systems, this translates into 10 billion complex multiplications per seconds, which requires a very high power consumption and/or silicon area. To reduce the burden of the DFE implementation, frequency domain (FD) techniques have been proposed, which enable performing so-called fast convolutions (or correlations), thereby significantly reducing the implementation complexity.\nConsidering the use of FD processing for the linear equalizer is quite straightforward because it is a direct application of the fast convolution or correlation. On the contrary, FD processing for the feedback part of a DFE has been rarely and not always convincingly addressed. The FD DFE presented in the papers “Blind Decision Feedback Equalization for Terrestrial Television Receivers” (M. Ghosh, Proc. IEEE, vol. 86, no. 10, pp. 2070-2081, October 1998) and “Overlap and Save Frequency Domain DFE for Throughput Efficient Single Carrier Transmission” (S. Tomasin, IEEE 16th Int'l Symp. Personal, Indoor and Mobile Radio Communications, September 2005) only have the feedforward part in the FD. The main reason is that the FD block processing does not lend itself easily to the feedback process. The FD approach in “Frequency Domain Feedforward Filter Combined DFE Structure in Single Carrier Systems over Time-varying Channels” (B. Liu et al., IEEE Trans. Cons. Electronics, vol. 54, no. 4, November 2008) targets mainly single-carrier (SC) systems with cyclic extension and is also limited to the FF part. In “Frequency-domain and Multirate Adaptive Filtering” (J. Shynk, IEEE Signal Proc. Magazine, vol. 9, no. 1, pp. 14-37, January 1992) both the FF and FB part, including the weight update, are done in the FD, but the FB is less performant than in the TD case because the new decisions obtained after the DFTs are not fed back until the next DFT is computed.\nRegarding notational conventions, normal Latin characters are used for time-domain signals (a) and tilde characters for frequency-domain signals (ã), vectors and matrices are denoted by a single and double under-bar, respectively, (a and A). The superscripts *, T and H denote the complex conjugate, the matrix transpose and complex conjugate transpose, respectively. The Hadamard (i.e. element-wise) product of vectors is denoted by the ⊙ operator. Any matrix in this description represented by a Q followed by a subscripted letter denotes a zero padding matrix. The operator F represents a FFT block operation.\nBefore a decision feedback equalization scheme as known in the art is described more in detail, a system model is introduced. The transmission is considered of a sequence of possibly complex symbols over a complex multipath channel and affected by additive white Gaussian noise. The discrete-time equivalent of the received signal x(t) can be expressed as:\n                              x          ⁡                      (            k            )                          =                                            ∑                              l                =                0                                            L                -                1                                      ⁢                                          s                ⁡                                  (                                      k                    -                    l                                    )                                            ⁢                              h                ⁡                                  (                  l                  )                                                              +                      n            ⁡                          (              k              )                                                          (        1        )            where s(l) is the sequence of the transmitted symbols and h(k) denotes the channel impulse response of length L. The received signal x(k) is fed to a DFE (see FIG. 1) that satisfies the following relationship\n                              y          ⁡                      (            k            )                          =                                            ∑                              i                =                0                            A                        ⁢                                                            a                  i                  *                                ⁡                                  (                  k                  )                                            ⁢                              x                ⁡                                  (                                      k                    +                                          d                      a                                        -                    i                                    )                                                              +                                    ∑                              i                =                1                            C                        ⁢                                                            c                  i                  *                                ⁡                                  (                  k                  )                                            ⁢                              d                ⁡                                  (                                      k                    -                    i                                    )                                                                                        (        2        )                                          d          ⁡                      (            k            )                          =                  f          ⁢                      {                          y              ⁡                              (                k                )                                      }                                              (        3        )            where ai(k) are the A coefficients of the feedforward filter, ci(k) are the C coefficients of the feedback filter, da is the delay of the feedforward filter and f{.} is the demodulation operator (slicer—in FIG. 1 named ‘decision device’). In many practical systems, the channel is not known at the receiver and it may be difficult to estimate the channel directly. Adaptive algorithms are then used to adapt the weights ai(k) and ci(k) so as to reduce the variance of the estimation error. Although many high performance algorithms exist, the least-mean square (LMS) algorithm is often used in practice for complexity reasons, especially when the filter lengths A and C are large. The LMS adaptation rule is as follows:e(k)=d(k)−y(k)  (4)ai(k+1)=ai(k)+2μax(k−i)e(k)  (5)ci(k+1)=ci(k)+2μcx(k−i)e(k)  (6)\nThe LMS algorithm used in combination with the proposed implementation of the DFE provides the simplest adaptation algorithm. However, its complexity increases as the channel length becomes larger. This is because the filtering operations in (2) involve convolutions of length A and C and the updates in (5) and (6) involve correlations of length A and C as well. Since the filter length must be approximately equal to or even greater than the channel length, the total complexity of the LMS-DFE is approximately equal to 4L possibly complex multiplications per output symbol, which may be challenging for channel lengths with hundreds of taps. It has been identified that the frequency domain (FD) processing is an interesting alternative to the time domain (TD) convolutions and correlations as in (2), (5) and (6). Indeed, it is well known that for a length Lx sequence x and length Ly sequence yconv(x,y)=IDFT(DFT(x)⊙DFT(y))  (7)corr(x,y)=IDFT(DFT(x)*⊙DFT(y))  (8)provided that the discrete Fourier transform (DFT) and inverse discrete Fourier transform (IDFT) are taken over a length larger than or equal to Lx+Ly−1. Zeros must be padded to x and y before taking the DFT to achieve this. This alternative implementation is attractive for large Lx and/or Ly if the fast direct and inverse Fourier transforms (FFT and IFFT) are used. When one of the sequences is infinite, it must be split into blocks and special techniques, known as overlap-and-save and overlap-and-add, must be used to recombine the blocks after (7) or (8) is applied.\nHence, there is a need for an adaptive feedback equalization scheme of reasonable complexity, which is suitable for applications where relatively long impulse responses are encountered."}
{"text": "The present invention relates to a novel system for applying for surgical clamps or clips which incorporates a novel disposable cartridge assembly detachably mounted on an instrument.\nApplicators for applying surgical clips or clamps which utilize or include a cartridge are known in the art, see for example U.S. Pat. Nos. 3,777,538; 3,232,089; and 2,968,041. The advantages of applicators of this type are enumerated at length in the specifications of the three patents noted above. The prior art, however, does not suggest a system for applying surgical clips which incorporates a disposable cartridge detachably mounted on an instrument in the unique manner of the present invention."}
{"text": "(1) Field of the Invention\nThe present invention relates to a tire shield device mounted between the tire and rim which resists penetration of the side wall of the tire by foreign objects such as glass, stones, wire, steel and the like which would otherwise penetrate the side wall. The tire shield device is also constructed and mounted so as to allow a vehicle to be operated at high speed without causing the tire or the shield to overheat or come apart.\n(2) Prior Art\nThe general prior art is shown by U.S. Pat. Nos. 3,811,488 to Duncan; 4,030,530 to Curtiss; 4,111,250 to Bauer; 4,235,271 to Olsen; 4,252,169 to Watts; and 4,319,618 to Suzuki. These patents show various devices for protecting the sidewall of off the road construction tires, none of which are high speed and travel an average of 5 mph. These later patents are more concerned with penetration protection for the sidewall of slow moving construction tires as used in mining, lumbering, scrap yards etc. U.S. Pat. Nos. 1,905,674 to Babbs; 1,875,067 to McKey and 2,017,891 to Brigqs show splash guards or skid protectors which are not designed for protecting the sidewall.\nSome of the prior art patents are concerned with providing sidewall protection which is integral with the tire. This results in a relatively expensive and heavy tire. In other instances the prior art is concerned with sidewall protectors which bolt to the surface of the rim which supports the tire. When the tire shield is connected to the rim, considerable weight is added to the tire and rim, it is inconvenient and time consuming to mount the shield and the holes in the shield for mounting create failure points for the shield.\nU.S. Pat. No. 3,187,797 to Fletcher et al describes a device for deflecting water from aircraft wheels. This device connects between the tire bead and the rim and is designed to extend based upon centrifugal force of the wheel and water pressure from under the tire. The device is not designed to resist penetration of the tire sidewall, since it extends only partially up the sidewall of the tire and it is intended only for aircraft use limited to landings or takeoffs.\nThe problem with the prior art tire shields is that they can come in contact with the sidewall of the tire in use and abrade the sidewall. Further, there can be a heat buildup between the shield and the tire where the shield is in contact with the sidewall at high speeds. Further still, the prior art shields tend to provide limited resistance to penetration because of their construction. This is why they are used for slow off the road vehicles.\nIn my prior U.S. Pat. Nos. 4,262,719, 4,263,074, 4,303,114 and 4,347,884, I described a method for forming a tire liner. The method of the present invention is different from the method of these prior applications because of the need to form a disc shaped shield."}
{"text": "Field of the Invention\nThe present invention relates to pneumatic tires, and in particular relates to a pneumatic tire being provided with a middle lug groove including a sipe portion formed between a protrusion on a first groove wall and a second groove wall.\nDescription of the Related Art\nJapanese Unexamined Patent Application Publication No. 2011-31773 discloses a pneumatic tire capable of improving wet performance and noise performance in good balance while ensuring steering stability.\nThe tire includes a tread portion provided with a crown main groove and a shoulder main groove to form a crown rib, a middle rib and a shoulder rib. The middle rib is provided with lug grooves extending from the shoulder main groove to improve drainage property of the middle rib. Furthermore, the middle rib has sufficient rigidity to ensure steering stability in high level since the lug grooves include axially inner ends that terminate within the middle rib. In addition, since the tire prevents that an air compressed between a lug groove and the ground flows into the crown main groove, the pipe resonance generated from the crown main groove, which is a main cause of road noise, can be reduced.\nHowever, since the above mentioned lug grooves include outer ends having a large width, it has been insufficient to prevent reduction of the rigidity of the middle rib. Furthermore, the air compressed by the lug grooves flows into the shoulder main groove to promote the pipe resonance. Therefore, noise reduction effect has also been insufficient.\nAs illustrated in FIGS. 9A and 9B, Japanese Unexamined Patent Application Publication No. 2013-233822 discloses a tire tread provided with a shoulder main groove (a) and a lug groove (b), wherein the lug groove includes a tie-bar (c) provided with a sipe (d). Unfortunately, when a depth (h) of the tie-bar (c) is shallow, sufficient drainage performance cannot be ensured since water would be blocked by the tie-bar. On the other hand, when the depth of the tie-bar (c) is deep, it may be difficult to ensure steering stability and noise performance of the tire. Accordingly, it has been difficult to improve steering stability, wet performance and noise performance at a high level."}
{"text": "(a) Field of the Invention\nThe invention concerns a process for preparing a plastic supported metallic sheet, obtained by metallization-plating, as well as the metallic sheet resulting from this process. More specifically, the invention is concerned with a process for the preparation of metallic films consisting of at least one metal which is deposited by metallization on a non-conductive plastic film, the product obtained being covered with another metal by electrochemical plating, these metallic films having a thickness within the range of 0.1 to 4 microns.\n(b) Description of Prior Art\nThe manufacture of thin metallic films is faced with difficulties when the thicknesses used are of the order of a micron. These difficulties signify high costs and leave only a limited choice of metals to be used. The three main industrial methods actually used are laminating, electro-forming on a mandrel and metallization. The minimum thickness obtained by the first two methods are of the order to 5 to 20 microns and the film produced is free. Metallization should be carried out on a plastic support and the thickness normally obtained on plastic does not exceed 0.1 microns by much.\nAmong the industrial processes presently available for the manufacture of thin films less than 15 microns, laminating is the more currently used. The latter consists in reducing the thickness of the sheet by passing it between two rollers on which a pressure is applied. However, few metals are laminated at these thicknesses (for example: Au, Al, Pb, brass, stainless steel) and, generally, the costs increase rapidly and the capacity of the processes decreases rapidly for thicknesses below 50 to 5 microns, by virtue of the number of passes required for laminating, the fact that the films having a low mechanical resistance are hard to handle, the appearance of perforations during the processes, as well as problems of reheatings and contamination with lubricants. Moreover, the minimum thicknesses which are commercially available (12 microns for stainless steel, 6-7 microns for aluminum) are still not acceptable and are penalizing in cost, material, weight or dead volume for many applications. The final product obtained by laminating is a coil of a rectangular metallic sheet.\nElectro-forming is a plating technique wherein the free and self-supported metallic film is obtained by electrochemical deposit, from an electrolytic solution, on an electronic conductive mandrel. The manufacturing costs are high (about 1$ US/ft.sup.2 for nickel) and the minimum available thicknesses are 6 microns for nickel and 10 microns for copper. This technique is applicable to certain materials only and, the necessity to handle and peel a free film after deposit, requires minimum thicknesses otherwise it would not be possible to peel the film from the mandrel without running the risk of damaging the film. Electro-formed copper films are produced in large quantity and are mainly used for the manufacture of printed circuits. The final product obtained by electro-forming a metallic sheet is a coil of rectangular shape.\nMetallization or vacuum deposit consists in evaporating a metal for depositing same on a support. The most current processes are: metallization by vaporization and assisted vaporizations such as cathodic pulverization and pulverization by beam of electrons. The process of vaporization is mostly used. It consists in heating the metal so that its vapor pressure be significantly higher than the pressure in the vacuum chamber. The quantity of metal that may be evaporated is dependent on the temperature used (1200.degree. C. for Al and copper) and the surface of the bath. The maximum thicknesses deposited are dependent on the process and the heat resistance of the substrate. The thicknesses of the metallizations on plastic materials rarely exceed 0.5 micron and are generally of the order of 0.01 micron. The thickness of the metallization is limited inter alia by overheating of the plastic. The deposits obtained on plastic are currently used in applications where the metallized thickness requires no thick deposits, i.e. &lt;0.1 micron, such as decorative metallizations, barrier against gas diffusion, current collectors of plastic condensers, partially reflecting coatings, etc. These metallizations are fragile, and, because of their small thickness, they are also not very conductive of electricity, generally of the order 0.5 to &gt;100 ohm/square.\nPlating on plastic is used, inter alia, in the automobile industry. The process is used in 85% of the cases to plate molded ABS articles. The first steps required before plating are: cleaning in an alkaline solution, rinsing, oxidation of the surface by means of a chromic based solution, neutralization to remove hexavalent chromium, rinsing, tin and palladium based catalytic bath, accelerated to remove tin and expose palladium, rinsing, autocatalytic plating with nickel or copper. After these steps, the articles are covered with about 0.12 to 0.75 micron of copper and may receive other plating.\nThe thicknesses which are proposed according to the present invention are not accessible to laminating, and metallization under vacuum enables at the most with certain heat resistant plastics to reach values lower than the range provided by the invention and this, at the price of a substantial slowing down of the process.\nContrary to the invention, the process of plating on plastic used in the prior art implies the use of a specific plastic and a number of chemical treatments. The process of plating on plastic requires apparatuses for controlling pollution, because of the use of metals such as hexavalent chromium. The raw material used in the process of plating on plastic has a rough surface because of the nature of the process according to which the plastic is oxidized and catalyst sites are produced for the autocatalytic deposit.\nThe invention aims at a process according to which the surface condition of the metallic precursor which is deposited under vacuum perfectly follows the plastic surface, which enables to obtain the desired surface.\nAnother object of the invention is to permit the treatment of only one face of the plastic film when this is desirable, while with the plating on plastic of the prior art all the surfaces of the articles are treated.\nAnother object of the invention is to enable an operation which is very regular ensuring a very precise control of the thicknesses, this being much more difficult with plating on plastic because of the number of steps and the notion of activation sites.\nAnother object of the invention resides in a process which is more simple than plating on plastic because of the limited number of operations to be carried out.\nAnother object of the invention consists in effecting the plating of the metal according to two methods: plating with imposed current and autocatalytic plating. Plating with imposed current enables speeds of plating which are much more rapid and the purity of the metal deposit is higher than with the autocatalytic plating.\nAnother object of the invention resides in the use of a conductive surface obtained by metallization under vacuum, which is very favorable to a metallic deposit which is rapid, uniform and chemically clean. No cleaning of the metallized surface is required before plating.\nAnother object of the invention consists in the production of sheets having a composite metallic structure; this possibility enables to optimize the properties of the product.\nAnother object of the invention is to select the last plated metal as a function of its resistance to corrosion, its weldability, its compatibility with a glue, etc. The underlying metal(s) are selected as a function of their conductivities, their mechanical or electrical properties, etc.\nAnother object of the invention is to control the shape and surface conductivity of the electrodeposited zones as a result of the utilization of well known techniques in galvanoplasty; this possibility has an advantage as compared to laminating. This advantage may be used to optimize the utilization of the materials.\nAnother object of the invention is to facilitate the handling in automatically operated apparatuses, of the metallic sheets produced, through the utilization of a plastic support.\nAnother object of the invention is to also enable to control the surface, thickness and shape of the parts where the electrochemical deposit will take place.\nAnother object of the invention is to enable to obtain in a same support, insulating zones, metallized zones, and zones having an electrochemical deposit of at least one metal, which is impossible with the processes of lamination, metallization and plating.\nAnother object of the invention is to additionally and optionally enable, the easy handling and the peeling of metallic coatings by transfer of adhesive materials with an appropriate control of the adhesion between the plastic support and the first deposit obtained by metallization under vacuum.\nAnother object of the invention resides in the utilization of metallic sheets for the manufacture of heating elements, thin metallic coatings to absorb or reflect light and/or heat metallic barriers to block the diffusion of gases in multi-layer wrappings and current collectors, such as in thin film batteries.\nAnother object of the invention is to use a metallic sheet as current collector of the electrodes of a thin film electrochemical generator.\nThe invention concerns a process for the manufacture of metallic sheets consisting of at least one metal supported on a non-conductive film of synthetic resin comprising the following steps:\nmetallization under vacuum, of at least one face of the plastic film so as to give a substrate including a metallic surface, said substrate being sufficiently electrically conductive to provide a uniform electrochemical deposit, PA1 deposit by electrochemical plating, of at least one metal on the metallized surface starting from an electrolytic solution, so as to give a thin metallic film having a metal thickness between 0.1 and 4 microns, the metallized substrate being selected so as to be compatible and facilitate the step of electrochemical plating, the thin metallic film obtained being supported by said plastic film. The metallization and the deposit by electrochemical plating may be carried out continuously.\nThe process of metallization under vacuum is preferably selected among simple or assisted thermochemical vaporization processes.\nWith respect to the assisted thermochemical vaporization, it is preferably carried out by cathodic polarization or with a beam of electrons.\nAccording to an alternative embodiment, the metallized surface is controlled by concealing or demetallization processes, for example by laser beam, chemical dissolving, chemical removal or flashing.\nThe insulating non-conductive plastic film is preferably made of a synthetic resin selected for its compatibility with the vacuum conditions and the solutions associated with the process of galvanoplasty, for example polypropylene, polyester, polysulfone, polyethylene, or polyimide.\nThe substrate includes a metallized surface and preferably consists of substances which are compatible with the deposit by electrochemical plating, selected among the following metals: Cu, Au, Ag, Fe, Ni, Cr, Zn, Mo or alloys thereof, particularly Cu, Au, Ag or alloys thereof, whose surface resistance varies for example between 0.1 and 10 ohm/square.\nAccording to another embodiment of the invention, metallization under vacuum is carried out so that the metallized surface covers only partially the surface of the insulating support film to leave non-conductive zones. A partial covering of the surface of the film by metallization of the insulating support film before the step of electrochemical deposit, may be obtained by the utilization of masks during metallization, by removal of the metallized deposit, for example by laser treatment, flashing, chemical or mechanical treatment.\nIt is also possible to cover the metallized surface with insulating masks to prevent local plating, which masks will, in certain cases, be removed thereafter. The metal which has been deposited by plating is for example selected from the following elements: Ni, Fe, Au, Ag or alloys thereof. Preferably, metallization is carried with copper and the metal deposited by plating comprises nickel. The metal deposited by plating is for example selected from the following elements: Ni, Fe, Au, Ag, Cr, Zn, Pd, Cd or alloys thereof.\nThe electrolytic plating solution of the autocatalytic type enables the utilization of a thinner metallization, than when an ordinary electrolytic solution is used. The metallic sheet may be covered with a plastic film, for example by laminating or gluing.\nAccording to another embodiment of the invention, the metallic sheet is transferred entirely or partially, from the plastic support towards another surface whose adhesion with respect to the last plated metal is higher than that between the metallization and the plastic support. The metallic sheet may also be chemically treated to provide it with a low emissivity in ultraviolet.\nThe metallic sheets according to the invention may be used to constitute heating elements which may be integrated in a flexible or rigid unit, current collectors for batteries, for example polymer electrolyte lithium batteries, wrappings with a very low permeability, light screens, flexible conducts, etc.\nThe capacity of the metals to block, reflect or absorb radiations is well known. In some cases, the use of a flexible metallic film is very desirable. The product according to the invention combines, for a plurality of applications, the advantages of a plastic film for handling, of the metallic surface for opacity. The film may be used in different types of blinds to block light.\nA modification of the surface of the sheet of nickel may, for example, produce a low emissivity of the ultra-violet rays. The use of a minimum thickness of nickel enables to preserve an excellent flexibility as compared to an electro-formed product. This product may be used to prevent systems for watching by night to be detected. It is also possible to coat the reverse side of the support film with an adhesive to make it possible to cover the surface to be protected.\nThe capacity of reflecting radiations may be used for utilization as vapor barrier. The facility of use of the plastic is then combined with the impermeability to gas, the reflectivity and the resistance to corrosion of the metallic coating such as nickel. The vapor barrier is therefore more efficient on a point of view of energy, duration and permeability.\nThe electrical resistance of the metallic film on plastic enables it to be used as a resistance to produce heat. This resistance may be used in ceilings for heating by radiation or as large surface heating element to heat liquids; the large heating surface prevents the formation of hot points which could deteriorate the content of the container.\nIt is preferable to select the plastic support as a function of the desired adhesion between the metallization and said support. To increase adhesiveness, the surface to be metallized may be treated by corona effect.\nMetallization is normally carried out under vacuum by utilizing a metal which promotes the later electrochemical deposit so as to induce a sufficient electrical resistance to permit the flow of plating current (preferably lower than 10 ohm/square). If the electrical resistance is too high, the plating current should be limited because of the danger of overheating, which would decrease by the same amount the maximum speed of plating. It is possible, by the utilization of masks, to obtain non-metallized surfaces on the plastic support, for example, to obtain a non-metallized border along the length of the film.\nA local covering of the metallization may be used to preserve non-plated metallized zones. These less conductive zones may be used as fuses, heating elements, or be later on dissolved.\nThe local demetallization may be carried out by removal with a laser beam, chemical solution, mechanical removal or flashing, when, for certain applications, the use of masks with the metallization appears to be difficult to realize.\nThe electrochemical deposit on the metallized film is carried out with at least one metal on the surfaces which are made conductive by metallization until obtaining a deposit of a thickness between about 0.02 and 4 microns. The conditions which are imposed during the electrochemical deposit are well known to one skilled in the art and are based on the electrolytic solution, the arrangement of the electrodes, the densities of current used which enable to optimize the density, conductivity and the surface condition of the electrochemical deposit. The plating speed which is possible with this process, and this without optimization, may reach 2 m/min when a thickness of 0.3 micron of nickel is deposited. Of course, during this step, it is possible to control the finish of the surface, the thickness and the resistivity of the consecutive electrochemical deposit(s), and this is made possible by varying the known parameters of the devices for electrochemical plating (current density, additives, covering, distance of the electrodes, etc.).\nThe demetallization of the plating deposit and/or of the metallization may be carried out by removal with a laser beam, mechanical removal or flashing when the demetallization is more practical than the use of a group of masks. Removal of the masks previously disposed may be carried out by any known means."}
{"text": "Retaining walls are commonly employed to retain highly positioned soil, such as soil forming a hill, to provide a usable level surface therebelow such as for playgrounds and yards, or to provide artificial contouring of the landscape which is aesthetically pleasant. Such walls have been made of concrete blocks having various configurations, the blocks generally being stacked one atop another against an earthen embankment with the wall formed by the blocks extending vertically or being formed with a setback. Setback is generally considered to be the distance in which one course of a wall extends beyond the front of the next highest course of the same wall. Concrete blocks have been used to create a wide variety of mortared and mortarless walls. Such blocks are often produced with a generally flat rectangular surface for placement onto the ground or other bearing foundation and for placement onto lower blocks in erecting the wall. Such blocks are also often further characterized by a frontal flat or decoratable surface and a flat planar top for receiving and bearing the next course of blocks forming the wall.\nIt is generally desired that retaining walls of the type described exhibit certain favorable characteristics, among which may be mentioned the ease with which the retaining wall can be assembled, the stability of the wall (that is, its ability to maintain structural integrity for long periods of time), and the ability of the wall to admit and disburse rainwater. Although retaining wall blocks commonly are supported vertically by resting upon each other, it is important that the blocks be restrained from moving outwardly from the earthen wall that they support.\nCurrent manufacturing techniques and the economics associated therewith limit the shapes, sizes, and materials that may be used to manufacture blocks that still provide the functions described above. In some instances, it would be preferred to make blocks in different shapes, sizes, and colors, and using different quality, types, and price of materials, and possibly in a centralized location which may be further from their point of use. It is desirable to both break through these boundaries and yet produce improved retaining wall blocks."}
{"text": "1. Field of Invention\nThe present invention relates to methods of forming an insulating layer and forming a device, and more generally to methods of forming a hydrophobic silicon dioxide (SiO2) layer and forming an organic thin film transistor (OTFT) including the hydrophobic SiO2 layer.\n2. Description of Related Art\nOTFTs have recently attracted much attention for their applications in radio-frequency identification (RFID) tags, large-area sensors and active matrix flat panel display (AMPD) backplane driving circuits.\nGenerally, a SiO2 layer is one of the most widely used gate insulating layers in OTFTs. When an OTFT is formed on a flexible substrate, the process temperature of the SiO2 gate insulating layer must be lowered for compatibility, otherwise the flexible substrate may be damaged. However, the SiO2 gate insulating layer formed at low temperature usually includes a significant amount of hydroxyl (OH) groups, which makes the surface of the SiO2 gate insulating layer more polar and hydrophilic. This OH-terminated SiO2 surface may cause poor crystalline growth of an organic semiconductor active layer (e.g. pentacene), and thus the performance of the OTFT is affected.\nIt is known that the OTFT performance can be greatly improved by modifying the SiO2 surface with a hydrophobic self-assembled monolayer (SAM), such as hexamethyldisilazane (HMDS) or octadecyltrichlorosilane (OTS). HMDS or OTS can form a hydrophobic methyl SAM (CH3-SAM) on the SiO2 surface, so as to reduce the SiO2 surface energy and improve the crystalline growth quality of the organic semiconductor active layer (e.g. pentacene). However, due to the surface modification treatment, the process becomes more complicated and time-consuming and the cost is increased."}
{"text": "Several publications and patent documents are referenced in this application to more fully describe the state of the art to which this invention pertains. The disclosure of each of these publications is incorporated by reference herein.\nThe p53 tumor suppressor protein is a nuclear phosphoprotein that functions in cell-cycle arrest, programmed cell death (apoptosis), inhibition of tumor growth, and preservation of genetic stability. The p53 gene is inactivated by mutation in over 60% of human tumor samples analyzed. This gene continues to hold distinction as the most frequently mutated gene in human cancer. The strong selection for mutation of p53 in cancer cells can be explained by the normal functions of the protein. p53 protein is normally expressed at very low levels in cells and is latent in activity. The protein is post translationally stabilized and activated as a transcription factor by a number of stimuli that are detrimental to the normal cell, notably genotoxic stress, hypoxia, and the mutational activation of oncogenes like c-myc and Ha-ras. The result of p53 induction in a cell is either growth arrest in the G1 phase of the cell cycle, or programmed cell death (apoptosis). The choice of fate depends upon a number of factors, including cell type, level of p53 induced, and environmental factors, such as the presence of cytokines or increased expression of bc12 family members. It is currently accepted that it is the ability of p53 to induce cell death, not G1 growth arrest, that underlies the powerful selection for mutation of the p53 gene in tumorigenesis.\nThe best-understood activity of p53 remains its ability to function as a sequence-specific transcriptional activator. Following a stimulus such as DNA damage, p53 protein is post-translationally stabilized and activated for DNA binding. This protein then binds in a sequence-specific manner to promoters or enhancers containing the consensus element 5′ Pu Pu Pu C A/T T/A G Py Py Py 3′. (SEQ ID NO: 12) p53 binding to these elements is believed to recruit the core transcriptional apparatus and thereby enhance the transcription of such (p53-response) genes.\nDespite the ability of p53 to transactivate pro-apoptotic genes like bax and IGF-BP3, in the past several years compelling indications have accumulated that p53 has a pro-apoptotic activity that is independent of its transactivation function. Recently it has been shown that physiological induction of wt p53 can lead to apoptosis in the absence of trans-activation of p53 response genes.\nIn addition to its well-characterized function as a transcriptional activator, p53 has a poorly-understood activity as a transcriptional repressor. Previous studies investigating this function of p53 were performed using transient over-expression assays where non-physiological levels of p53 and candidate promoters are introduced into cells. In this type of setting, a large number of promoters are repressed by p53. In no case, however, have these endogenous genes been found to be repressed by p53. Therefore, in these types of assays, transcriptional repression by p53 is non-physiological, and impossible to distinguish from transcriptional “squelching”. During transcriptional “squelching”, a powerful transactivator like p53 represses transcription of other promoters non-specifically by more efficiently recruiting basal transcription factors, such as TBP-associated factors (TAFs).\nTo determine whether p53 can function as a sequence-specific transcriptional repressor, differences in gene expression in cells containing an inducible p53 protein have been analyzed (Murphy et al., Genes & Devel. 10:2971–2980, 1996). Several genes were identified that exhibit decreased expression following wt p53 induction in cell lines containing a well-characterized temperature-sensitive mutant of p53. One of these genes encodes the microtubule-associated protein Map4. It was found that p53 induction leads to repression of Map4 at the level of transcriptional initiation (Murphy et al., Genes & Devel. 10:2971–2980, 1996). This transcriptional repression is manifested at both the RNA and protein levels, and approaches a 90% reduction at the RNA level following p53 induction for 24 hours. Further, physiological induction of p53, by DNA damaging agents or ionizing radiation, leads to significant repression of Map4 in both normal and tumor cells containing wt p53. In cells with mutant p53, Map4 levels are unaffected by such treatments (Murphy et al., Genes & Devel. 10:2971–2980, 1996). This and other data indicate that independent of it's cell cycle effects, wt p53 represses the expression of Map4. However, the precise mechanism by which p53 represses Map4 expression remains to be elucidated.\nThe Map4 gene is transcriptionally repressed following p53 induction in cells that undergo either G1 arrest or apoptosis in response to p53. Significantly, it has been found that removing Map4 from transcriptional repression by p53 significantly interferes with p53 dependent apoptosis (Murphy et al., Genes & Devel. 10:2971–2980, 1996), thereby directly placing Map4 in a pathway influential in p53-mediated programmed cell death.\nFollowing the cloning of Map4 as the first endogenous gene repressed by physiological induction of wt p53, several other genes have been cloned that exhibit decreased expression following p53 induction. These genes include those encoding DNA topoisomerase IIα, wee1, DP-1, presenilin 1, and others. It is clear from these studies that there will exist other candidate p53-repressed genes. Significantly, repression of at least one of these genes, presenilin 1, has been demonstrated to contribute to the progression of apoptosis. Further, the combined data from several sources suggest a perfect correlation between the ability of p53 to induce apoptosis, and its ability to repress the expression of genes like Map4.\nIt is clear from the foregoing that the expression repression function of p53 is a significant part of p53's central role in controlling cell growth and death. In order to utilize this function in a practical way, more information is needed relative to the physical interactions among p53 and its targets for repression in cells."}
{"text": "The present invention relates to a method of fabricating a scouring pad for a cleaning equipment, and more particularly to an efficient scouring pad fabrication method which doubles the productivity.\nConventionally, the fabrication of a scouring pad comprises the step of pouring a sponge dope 7 in a mold 9, the step of covering a scouring pad 8 on the sponge dope 7, and the step of heating the mold 9, enabling the sponge dope 7 to be foamed into a sponge 71. After foaming, the sponge 71 and the scouring pad 8 are formed integral with each other. This scouring pad fabrication method is less efficient."}
{"text": "This invention relates in general to a method and apparatus for pumping a cryogenic fluid from a storage tank. The apparatus comprises a reciprocating pump and the method comprises controlling pump flow rate and vapor pressure within the storage tank by controlling the proportion off cryogenic liquid and vapor supplied to the pump during the induction stroke.\nCryogenic fluids are defined as liquids that boil at temperatures of less than about 200xc2x0 Kelvin at atmospheric pressure, such as hydrogen, helium, nitrogen, oxygen, natural gas or methane.\nFor containing cryogenic fluids, vacuum insulated storage tanks are known. For example, liquefied natural gas (LNG) at pressures of between about 15 and 200 psig (about 204 and 1580 kPa) can be stored at temperatures of between about 120xc2x0 K and 158xc2x0 K in vacuum insulated storage tanks.\nA problem with known storage tanks is that heat leaks can cause vaporization of some of the cryogenic fluid within the storage tank and this reduces the time that cryogenic fluids can be held within such storage tanks. To prevent the vapor pressure from rising to undesirable pressures, cryogenic storage tanks are normally equipped with a pressure relief valve. When the vapor pressure rises to above the set point for the relief valve, the storage tank is vented. There is a need for a system that reduces the need for venting, since it may be undesirable to release some cryogenic fluids into the atmosphere and since venting is wasteful of cryogenic fluid.\nSome cryogenic fluids such as hydrogen, natural gas, and methane are usable as fuels in internal combustion engines. In some engines, improved efficiency and emissions can be achieved if the fuel is injected directly into the cylinders under high pressure at the end of the compression stroke of the piston. The fuel pressure needed to inject fuel directly into the engine cylinder in this manner can be 3000 psig (about 23,700 kPa), or higher, depending upon the engine design. Accordingly, the cryogenic fuel cannot be delivered directly from a conventional storage tank and an apparatus is needed for delivering a cryogenic fluid to the engine at such high pressures. A pump is required to boost the pressure from storage pressure to injection pressure and to remove vapor from the storage tank to reduce the need for venting.\nU.S. Pat. No. 5,411,374, and its two divisional patents, U.S. Pat. Nos. 5,477,690 and 5,551,488, disclose embodiments of a cryogenic fluid pump system and method of pumping cryogenic fluid. In one embodiment the disclosed double-acting piston pump may be employed as a mobile vehicle fuel pump. In this embodiment, the pump is employed to remove both cryogenic vapor and liquid from the tank in a manner whereby only liquid is removed when the pressure in the surge tank is low and vapor starts to be removed when pressure in the surge tank is sufficiently high for engine demand and the vapor pressure in the vehicle tank is higher than the set point. The cryogenic liquid and vapor are supplied from a storage tank through respective conduits communicating between the tank and the pump inlet. A liquid control valve is associated with the liquid supply conduit and a vapor control valve is associated with the vapor supply conduit. The liquid and vapor control valves are controlled in response to fuel demand and the vapor pressure measured within the cryogenic storage tank.\nCo-owned U.S. Pat. No. 5,884,488, which is hereby incorporated by reference herein in its entirety, discloses a high-pressure fuel supply system for supplying cryogenic fluid from a storage tank to an engine. The \"\"488 patent discloses, among other things, a multi-stage LNG pump that is capable of pumping liquid or a mixture of liquid and vapor. A metering valve is adjustable to control the amount of vapor drawn into the pump suction. In another embodiment, an orifice is provided in the vapor supply line for regulating the amount of vapor induced into the sump for the LNG pump. The technique disclosed herein permits increased holding times in the storage tank by providing a method and apparatus for removing vapor from the storage tank.\nIn the present method, cryogenic liquid and vapor is pumped from a storage tank with a reciprocating piston pump. The method comprises:\n(a) In an induction stroke,\nretracting a piston within the reciprocating pump and drawing cryogenic fluid from the storage tank into a piston chamber associated with the piston;\ncontrolling flow rate through the pump by controlling the proportion of liquid and vapor supplied to the pump by supplying substantially only vapor during a selected portion of the induction stroke; and\n(b) in a compression stroke, compressing and condensing any vapor and compressing any liquid within the piston chamber, and discharging compressed cryogenic fluid from the pump.\nIn a preferred method, flow rate through the pump is controlled to maintain pressure within a predetermined range at a point downstream from the pump. For example, the point downstream from the pump may be in an accumulator vessel, in a pipe, or in a manifold of a fuel system leading to an engine.\nThe method may further comprise monitoring vapor pressure within the storage tank and further controlling the proportion of vapor and liquid supplied to the pump to maintain vapor pressure within the storage tank below a predetermined value. For example, by changing pump speed, a constant flow rate may be maintained, while changing the proportion of liquid and vapor supplied to the pump. Similarly, when pressure downstream from said pump is within the desired predetermined range, the proportion of vapor supplied to the pump may be increased to reduce vapor pressure within the storage more quickly.\nThe proportion of liquid and vapor supplied to the pump during the induction stroke may be controlled by first supplying liquid until the piston reaches a position during the induction stroke that corresponds to a desired proportion of liquid and then supplying substantially only vapor to fill the piston chamber until the induction stroke is complete.\nIn a preferred embodiment, for each pump cycle, the minimum flow rate pumpable through the pump is determined by the minimum proportion of liquid that is needed during the compression stroke to allow condensation of the vapor within the piston chamber.\nA liquefied gas occupies much less space than the same fluid in the gaseous state, so a storage space advantage may be realized by applications that use cryogenic systems to supply a gas. For high-pressure applications a cryogenic pump may be employed. After the liquefied gas is discharged from a cryogenic pump, the fluid may be directed to a heater for transforming it into a gas.\nIn one embodiment of the method, the desired proportion of liquid, measured by volume, is constant in each pump cycle. To achieve a constant proportion of liquid, vapor is supplied to the pump during a predetermined portion of the induction stroke. For example, liquid may be supplied to the pump initially from the beginning of the induction stroke and whenever the piston reaches a predetermined position, vapor is then supplied to the pump for the remainder of the induction stroke. The same result could be achieved by supplying substantially only vapor to the pump during any predetermined constant portion of the induction stroke, and substantially only liquid during the rest of the induction stroke.\nWhen the cryogenic fluid is a combustible fuel, the present method may be employed to supply fuel to an engine.\nIn one embodiment, the supply of vapor to the piston chamber during the induction stroke is controlled by operating an automatically actuated valve associated with a vapor supply pipe that connects an ullage space of the tank with the pump. The method comprises opening the valve to supply substantially only vapor to the pump and closing the valve to supply substantially only liquid. The flow rate through the pump is controlled by controlling when the valve is opened with reference to the position of the piston, and flow rate is increasable by opening the valve for a smaller portion of the induction stroke. The position of the pump piston is determined by a sensor that sends an electronic signal to an electronic controller. The sensor may comprise a linear position transducer associated with the piston. Suitable means for automatically actuating the valve are well known. For example, the actuator may be electronic, mechanical, pneumatic, hydraulic, or a combination these. A mechanical actuator may be set to automatically actuate the valve for a constant portion of the induction stroke.\nIn a preferred embodiment, the valve actuator is electronically controlled and the proportion of liquid and vapor supplied to the pump is variable from one induction stroke to the next. For example, an electronic controller may be employed to open and close a solenoid actuated valve for directing vapor to the pump and achieving a desired pump flow rate. By supplying vapor from the ullage space of the storage tank to the pump, vapor pressure within the storage tank is reduced.\nAn advantage of the present technique is that a metering valve or orifice is not required to control the amount of vapor that flows through the vapor supply pipe. Instead, according to the present method, the proportion of vapor may be controlled in each individual induction stroke.\nIn a preferred method, a linear hydraulic motor drives the pump. A linear hydraulic motor is preferred compared to a mechanical crankshaft drive since a linear hydraulic motor can be used to operate the pump at a constant speed and this reduces pressure pulses in the discharge pipe. When the method is employed for supplying fuel to an engine, mechanical energy from the engine may be efficiently used for powering a hydraulic pump for the hydraulic motor.\nWhen a linear hydraulic motor drives the pump, the position of the pump piston may be determined by monitoring the hydraulic motor. In another embodiment, the position of the pump piston is determined by monitoring a reference point associated with the piston rod disposed between the pump piston and the linear hydraulic motor.\nWhen the method employs a single stage pump, at a given pump speed, the apparatus can be controlled to operate at a maximum flow rate by supplying only liquid to the pump during the induction stroke. When the pump is equipped with an inducer, an amount of vapor may still be supplied to the pump when the pump operates at a maximum flow rate because the vapor is condensed in the inducer. With an inducer, for each cycle, maximum flow rate is achievable by supplying a proportion of liquid and vapor to the inducer such that all of the vapor supplied to the inducer is condensable by the inducer and liquid discharged from the inducer fills the pump piston chamber.\nIn another embodiment, the proportion of liquid and vapor supplied to the pump may be controlled by controlling the flow rate of the liquid supplied to the pump. For example, when vapor is not being supplied to the pump a flow control valve associated with the liquid supply pipe may be operated to control the flow rate of liquid flowing from the storage tank to the pump. Accordingly, for a pump that is configured to supply vapor to the pump for a constant portion of the induction stroke, the proportion of liquid and vapor supplied to the pump is controllable by controlling the flow rate of the liquid supplied to the pump.\nIn addition to controlling flow rate by controlling the proportion of liquid and vapor supplied to the pump, flow rate through the pump may be further influenced by employing a variable displacement pump or by changing pump speed. For example, when the pump is driven by a hydraulic motor, a variable speed controller can be used to change pump speed. In arrangements where the hydraulic pump or the cryogenic pump itself is driven by an engine that is supplied with fuel by the cryogenic pump, engine speed generally correlates to fuel demand and the pump speed can be controlled to automatically increase with increased engine speed. However, in this arrangement, a hydraulic motor with a hydraulic pump driven by the engine has an advantage over a cryogenic pump directly driven by the engine, because the hydraulic motor permits the pump speed to be controlled to reduce pressure pulses in the discharge pipe.\nWhen a variable displacement cryogenic pump is employed, flow rate through the pump may be further controlled by changing pump displacement, for example, by limiting the stroke when a lower flow rate is desired. Persons skilled in the technology involved here will understand that many methods of controlling flow rate through the pump may be combined with the disclosed method of controlling flow rate by controlling the proportion of cryogenic vapor and liquid supplied to the pump.\nA specific preferred method of pumping a cryogenic fluid from a storage tank with a reciprocating piston pump comprises:\n(a) in an induction stroke,\nretracting a piston within the reciprocating pump and drawing cryogenic fluid from the storage tank into a piston chamber associated with the piston;\nsupplying vapor from the storage tank to the pump through a vapor supply pipe when a valve associated with the vapor supply pipe is open;\nsupplying cryogenic liquid from the storage tank to the pump through a liquid supply pipe when the valve is closed; and\nreducing vapor pressure within the storage tank and controlling pump flow rate by controlling the timing for opening the valve during the induction stroke; and\n(b) in a compression stroke,\nreversing the direction of the piston to compress and condense vapor and compress the cryogenic liquid within the piston chamber; and\ndischarging compressed cryogenic fluid from the pump.\nWhen the pump induces liquid at the beginning of the next induction stroke, the valve associated with the vapor supply pipe is closed prior to the next induction stroke. The valve may be closed upon completion of the compression stroke or at any time during the compression stroke. Obviously, when the vapor is supplied at the beginning or during the middle of the induction stroke the valve is closed prior to the end of the induction stroke.\nThe present technique is further directed to an apparatus for carrying out the method of pumping a cryogenic fluid from a storage tank and reducing vapor pressure within the storage tank. In a preferred embodiment, the apparatus comprises:\n(a) a reciprocating pump for pumping the cryogenic fluid supplied from the storage tank;\n(b) a liquid supply pipe that fluidly connects the storage tank to an inlet of the pump;\n(c) a vapor supply pipe that fluidly connects an ullage space within the storage tank to the inlet;\n(d) an automatically actuated valve associated with the vapor supply pipe, the valve being operable between a closed and an open position for allowing vapor to flow through the vapor supply pipe when the valve is in the open position; and\n(e) a controller for determining when to open the valve during an induction stroke of the pump, the controller making such determination to achieve a desired flow rate.\nThe apparatus may further comprise a position sensor for determining the position of a piston of the pump. The position sensor communicates with the controller so that the controller opens the valve when the piston is in a position that corresponds to the desired proportion of liquid for the induction stroke. In a preferred arrangement, the position sensor comprises a linear position transducer associated with the piston.\nThe reciprocating pump may further comprise an inducer. The inducer is fluidly disposed between the storage tank and the reciprocating pump. The inducer comprises an inlet for receiving cryogenic fluid from the storage tank, an inducer piston that is reciprocable within an inducer piston chamber for compressing and condensing cryogenic vapor and compressing cryogenic liquid, and an outlet for discharging the compressed cryogenic fluid. The cryogenic fluid compressed by the inducer is then supplied to the inlet of the pump for further compression of the cryogenic fluid.\nIn a preferred arrangement of the inducer, the inducer piston divides the inducer piston chamber into two sub-chambers so that the inducer operates with two stages. Cryogenic liquid is transferred from the first piston chamber to the pump piston chamber through a one-way flow conduit, which is typically a check valve. A pressure-actuated valve allows cryogenic fluid to flow from the inducer\"\"s second stage to the first stage when pressure within the second stage exceeds a predetermined value. That is, during the compression stroke of the second stage, cryogenic liquid is transferred from the second stage sub-chamber to the pump piston chamber, and when the pump piston chamber is filled, the pressure within the second stage sub-chamber rises until the pressure actuated valve opens to return the excess fluid to the inducer\"\"s first stage sub-chamber. Such a two-stage inducer configuration allows excess cryogenic fluid to be recycled within the inducer instead of being returned to the storage tank.\nCryogenic pumps comprising inducers are described in more detail and illustrated in co-owned U.S. Pat. No. 5,884,488, which has been incorporated herein by reference in its entirety. The pump piston chamber is preferably volumetrically smaller than the inducer piston chamber. More particularly, the inducer piston chamber preferably has a volume that is at least two times larger than the volume of the pump piston chamber, and in a preferred embodiment, the inducer piston chamber has a volume that is between about four and seven times larger than the volume of the pump piston chamber."}
{"text": "The present invention relates to controlling temperatures of a cooktop and more particularly to limiting the temperature of a glass-ceramic cooktop such that a maximum temperature of the cooktop reaches about a predetermined maximum temperature.\nThe new trend in electronically controlled cooktops and/or ranges, typically, includes a cooktop surface composed of a glass-ceramic material that is positioned above one or more radiant heating elements. The electronically controlled cooktop includes various user controls that are operated by a user to adjust the amount of heat and, ultimately, the temperature desired for cooking. The radiant heating elements can be powered by electricity, natural gas or other sources. Typically, the radiant heating element and the user controls are connected to a controller that controls the amount of heat supplied to the cooktop. The electronically controlled cooktop also includes temperature and other sensor that are connected to the controller to aid in controlling the heat supplied by the radiant heating source. The temperature and other sensors can also be used in conjunction with the controller to detect certain conditions that can arise during operation of the cooktop.\nTo increase the service life of the glass-ceramic cooktop, the temperature of the glass-ceramic cooktop should not exceed a predetermined maximum temperature for extended periods of time. Typically, the temperature of the glass-ceramic cooktop is affected by several factors during operation. The power level selected by the user using a user control interface is an important factor affecting the temperature of the glass-ceramic cooktop. Generally, the user power selection results in a temperature of the glass-ceramic cooktop that is below the predetermined maximum temperature. However, the user power selection along with other dynamics and/or factors involved with the operation of the glass-ceramic cooktop can cause the temperature of the glass-ceramic cooktop to rise above the predetermined maximum temperature. As mentioned above, exposure to a temperature above the predetermined maximum temperature for extended periods of time may reduce the service life of the glass-ceramic cooktop.\nThe other dynamics and/or factors that affect the glass-ceramic cooktop temperature include, such as, for example, the type of cooking utensil used, the thermal conductivity of the cooking utensil, the contents of the cooking utensil and user manipulation of all of these dynamics and/or factors. The type of cooking utensil affects the temperature of the glass-ceramic cooktop because the size and manufacturer of the cooking utensil can affect the thermal conductivity of the cooking utensil. The thermal conductivity of the cooking utensil relates to the ability of the cooking utensil to transfer heat from the radiant heating element to the contents of the cooking utensil. In addition, factors, such as, for example, the material composition of the cooking utensil, the thickness of the cooking utensil, the flatness and/or warping of the bottom of the cooking utensil also affect the thermal conductivity of the cooking utensil. These factors affect the contact interface area between the cooking utensil and the glass-ceramic cooktop, and the dynamics of this interface area are important factors in determining the thermal conductivity of the cooking utensil. In addition, the contents of the cooking utensil or lack thereof can affect the temperature of the glass-ceramic cooktop because, for example, a cooking utensil that has boiled dry can cause the glass-ceramic cooktop temperature to increase. Additionally, the user can alter or manipulate some or all of these factors/dynamics during the cooking process by, for example, adding more contents to the cooking utensil and/or moving or rocking the cooking utensil on the glass-ceramic cooktop surface during the cooking process.\nAs stated above, various factors and/or dynamics can cause the temperature of the glass-ceramic cooktop to exceed the predetermined maximum temperature, and therefore, potentially reduce the service life of the glass-ceramic cooktop. Therefore, there is a desire to control the power supplied to the radiant heating element such that a maximum temperature of the glass-ceramic cooktop reaches about the predetermined maximum temperature and/or does not exceed the predetermined maximum temperature within accepted tolerance limits."}
{"text": "This invention relates to the art of fluid pressure regulating valves wherein a valve element is diaphragm actuated to control the flow of fluid through the valve and, more particularly to an improved valve element operating arrangement for a diaphragm-operated pressure regulating valve.\nIt is well known to provide fluid pressure-regulators which operate in accordance with the pressure of fluid acting against a flexible diaphragm between valve housing components to provide a pressure chamber on one side of the diaphragm and a chamber vented to atmosphere on the other. The pressure chamber includes an inlet orifice or port for fluid under pressure from a source and an outlet passage for the flow of gas to a point of use such as, for example, a gas-operated cooking grill. A valve unit controls the flow of fluid through the inlet orifice to the outlet opening and, most often, the valve unit comprises lever arm pivotally mounted in the pressure chamber and having an end connected to the diaphragm for the later to pivot the arm about its axis during operation of the valve. The other end of the lever arm carries a resilient sealing member or valve element which, in response to pivotal displacement of the lever arm, engages and disengages the inlet orifice to respectfully close and open the latter to the flow of fluid into the pressure chamber from the source. Regulating valves of the foregoing character also include a relief valve arrangement by which excessive back pressure into the pressure chamber from the point of use is vented to atmosphere.\nThe problems encountered in connection with pressure regulators employing the lever arm concept include an undesirably large housing or envelope and therefore a higher than desired cost of manufacturing. In particular in this respect, the length of the lever arm required to obtain the necessary mechanical advantage with respect to controlling the flow of liquid at a given supply pressure requires a large housing or envelope, and the size of the latter increases as the supply pressure increases requiring a larger capacity inlet orifice and liquid flow and a high lever ratio to assure appropriate control of the flow. It follows that any increase in the size of the valve housing increase the cost-thereof as well as the overall cost of the regulating valve. A further problem encountered in connection with lever operated valve elements is performance limitations resulting from the opening and closing arc of the sealing element relative to the inlet orifice and which arc is determined by the length of the lever arm between the pivot axis thereof and the valve element thereon. More particularly in this respect, the displacement of the sealing member or valve element along an arcuate path relative to the inlet orifice makes it extremely difficult to obtain and to then maintain proper seating of the sealing element relative to the orifice. Yet a further problem encountered in connection with pressure regulating valves using the lever concept results from turbulent flow conditions resulting from the fluid flow through the inlet orifice to the pressure chamber and which turbulence causes the lever arm to vibrate and/or oscillate during valve operation causing fluctuating gas flow to the point of use.\nIn accordance with the present invention, an improved valve element operating arrangement is provided for a diaphragm-operated pressure regulating valve which eliminates the use of a lever and, in connection with a pressure regulating valve of the same capacity using the lever concept, considerably reduces the size of the valve housing and thus the cost thereof as well as the overall cost of the regulating valve while increasing performance characteristics relative to lever operated valves. More particularly in accordance with the invention, the foregoing attributes are achieved through the use of a valve element which is coaxial with and axially displaceable relative to the inlet orifice and the opening and closing displacement of which is controlled by a ball and cam actuating unit which is coaxial with the inlet orifice and operates in response to displacements of the diaphragm. The ball and cam unit provides the necessary mechanical advantage which in turn provides the operating force against the valve element and the fluid supply pressure for closing the inlet orifice. By having all of the component parts of the actuating unit coaxial with and linearly displaceable relative to the axis of the inlet orifice, the housing or envelope for the valve is considerably reduced in size, and thus cost, relative to a valve of the same capacity employing the lever concept. Moreover, linear displacement of the valve element relative to the inlet orifice advantageously provides for a resilient sealing member on the valve element to engage the orifice transverse to the axis thereof, thus optimizing the obtaining and maintaining of a proper seating of the sealing element against the inlet orifice.\nIn accordance with another aspect of the invention, the flow of fluid from the inlet orifice to the outlet passage of the valve is through trim orifices, each of which is slightly smaller in diameter than the inlet orifice and which trim orifices function together with the inlet orifice to promote uniform flow, to eliminate vibration and to provide lockout pressure requirements. More particularly in this respect, the trim orifices create a positive pressure or force against the ball and cam components of the actuating unit to eliminate clearances therebetween and thus undesired relative movement therebetween. By reducing lockout pressure, an increase in diaphragm bias spring loading is obtained which results in increased regulated output flow, and this further contributes to a smaller sized regulator and a higher flow rate.\nIt is accordingly an outstanding object of the present invention to provide an improved diaphragm-operated valve element actuating arrangement which enables decreasing the size and cost of a diaphragm-actuated fluid flow pressure regulating valve in comparison to such valves employing a diaphragm actuated pivotal lever for displacing the valve element relative to the inlet orifice of the valve.\nAnother object is to provide a diaphragm actuated pressure regulating valve in which the valve element is coaxial with and axially displaceable relative to the inlet orifice of the valve and wherein the valve element actuating unit is coaxial with the valve element and operated by displacement of the diaphragm.\nA further object is the provision of a diaphragm actuated pressure regulating valve of the foregoing character which is structurally simple, of minimal size with respect to a given flow capacity, and which costs less to manufacture than diaphragm actuated fluid flow pressure regulators heretofore available.\nYet a further object is the provision of a diaphragm actuated fluid flow pressure regulating valve of the foregoing character in which the diaphragm operated valve element actuating unit is comprised of ball and cam components operatively interconnected with the valve element for controlling displacement thereof in response to displacements of the diaphragm."}
{"text": "1. Field\nThe Present Disclosure Generally Relates to the application of film over a substrate, and deals more particularly with a method and apparatus for applying and trimming a film material along a surface of an elongate structural member, such as a composite stringer.\n2. Background\nComposite and other structures may be bonded together using a structural film adhesive material which may include a carrier such as a scrim for controlling bondline thickness. For example, in the aircraft industry, a composite stringer may be joined to another composite part such as a skin by co-bonding, or in a secondary bonding operation using a structural film adhesive material that is applied over the surface of a flange on the stringer. The stringer flange with the applied film adhesive is placed against the skin and the film adhesive is cured using thermal or other techniques.\nPreparing long composite structural members such as aircraft stringers for an adhesive bonding operation is a multi-stage process that is currently time consuming and labor intensive. A film adhesive material is manually rolled out to the length of the stringer, which may be 90 feet or more in the case of aircraft applications, and a protective poly backer must then be removed from the film. The entire length of the film material is then flipped over by a team of workers in order to place the backside of the film adhesive on a base flange of the stringer. Workers then manually trim excess film material away from the edges of the base flange. This hand trimming technique may be tedious, time-consuming and presents the risk of damage to edges of the base flange, particularly where the flange is beveled. During hand trimming, foreign object debris may be produced which, if adhered to the stringer, may affect the desirability of the bond. After trimming, the film material is smoothed by hand following which the stringer is flipped over and vacuum bag processed in order to compact the film material against the stringer base flange.\nAccordingly, there is a need for a method and apparatus for applying film material to long structural members that reduce dependence on manual labor, and which improve process flow time. There is also a need for a method and apparatus for applying film material that provide rapid trimming of excess material to match structural details, and which eliminate the need for vacuum bag compaction."}
{"text": "During astronomical observations, amateur astronomers usually use one of such processes: (1) Determining the position of celestial bodies by consulting astronomical almanac, star maps, etc.; adjusting the telescope, looking for the celestial body for observing, after locating the celestial body, setting the photographic equipment to take photos. During the whole observation process, there are a lot of inconvenience, such as, one must consult the relevant information to determine the positions of celestial bodies before each observation as the celestial bodies are constantly moving in the sky, but the ephemeris is heavy and non-intuitive, the simple star chart can only used to observe at specified times due to the moving of the celestial bodies; to adjust the telescope manually is very complicated; and it's not convenient to take photos through the telescope. (2) Using a guiding software and an electric equatorial to achieve automatic star tracking. This method requires calibrating the position of the telescope first, the guiding process is complicated to operate, the interactive interface is unfriendly, and one should have much astronomical background knowledge to complete the operation.\nThe methods described above all need manual focus to get a good observation effect, it is trouble to operate when tracking celestial bodies or taking a long time exposure for photograph, and there is a high threshold for beginners so it will hinder the spread of astronomical observations.\nMeanwhile, when taking the photos of the celestial bodies, a dedicated interface needs to be used to connect the camera to the eyepiece; it is not convenient for immediate viewing and network sharing.\nIn the prior art, a Chinese patent application with a publication number CN1808207A discloses an automatic finder control device for a portable astronomical telescope, using a main controller with microprocessor and the like, an equatorial drive controllers, GPS and DC motor, to control the telescope for pointing and tracking the target object. The device structure is complex, the main controller obtains data of stars stored in FLASH, using the information of time and location obtained from GPS, to deal with the calculations such as annual parallax calculation, solar parallax calculation, aberration calculation, calculation of precession, nutation calculation, and atmospheric refraction computing, the control process is complicated, and the device can not achieve automatically focus and photograph of celestial bodies.\nA China Patent with a publication number CN101017240A discloses a video CCD auto guider method for an astronomical telescope, proposes to use the CCD for telescope auto guiding. However, in this method, the computer obtains video images by a video capture card, which means the telescope connects with the computer via wired connection, so it is typically used for the observatory and other fixed occasions, not suitable for using in mobile occasions by amateur astronomers. At the same time, this method also does not give an automatic focus method.\nFor ease of movement and control, a Chinese Utility Model with a publication number CN202334667U discloses an astronomical telescope automatic star finding and tracking device based on Android mobile phone control, using the Android phone to select objects to be observed on the touch screen, connected by WiFi to transfer a set of drive parameters to a control module, the control module controls the servo motor to rotate the telescope via the drive module to achieve star finding and tracking automatically. Although this means achieves a wireless control, and is easy to carry and use, but the device does not have the function of autofocus or the photograph of celestial bodies. If transferring the images via the WiFi in normal ways, the image data need for auto finding and auto focus is too large to transfer, it is difficult to achieve real-time image transmission and detection control.\nIn addition, the prior art lacks suitable auto focus method for portable astronomical observation system. As the distance of the celestial bodies is very large in astronomical observations, the focus method is different with such in conventional photography. The literature “An astronomical telescope automatic focusing system based on image definition evaluation” published in “Astronomical Research & Technology” 2008, 3, disclosed an automatic focusing method for telescope based on image definition evaluation. However, this method is only applicable to fixed stars, using only the near roundness and size of the stars as the basis for focusing, the relationship between the focusing function and the two parameters is not described, that is, the specific calculation method for focusing function is not given, and when using the method in focus searching, it is easy to fall into local maxima without finding the global optimum.\nTherefore, a new automated astronomical observation system should be provided to accommodate the needs of amateur astronomers."}
{"text": "1. Technical Field\nThe present invention relates to an array antenna having a structure in which a plurality of antenna elements is arranged in a predetermined fashion. More particularly, the present invention relates to a series-fed phased array antenna which can electronically steer an antenna beam through variation in frequency.\n2. Description of the Related Art\nGenerally, in a radio communication system, an antenna is used to send and receive signals and the length of an antenna depends on the frequency to be transceived. Such antennas have developed into various forms as technology evolves. Recently, many studies have been made to develop a method of using multiple antennas.\nParticularly, array antennas having a structure in which multiple antennas are arranged in a predetermined fashion are widely used. Such an array antenna is a device that uses a large number of radiating elements to acquire a narrow beam width and transmit a signal. When using the array antenna at broadband frequency, beam direction, efficiency, unit cost, and the like may vary depending on the type of feed network.\nGenerally, when a parallel feed network is used, the beam direction is fixed even when the frequency to be transmitted changes. However, as the overall length of the feed network increases, signal loss along a transmission line increases and transmission efficiency drops when the transmission line is constructed using a dielectric substrate. In addition, when a transmission line of a parallel feed network is constructed using a waveguide, the feed network becomes complicated, thereby making it difficult to manufacture the network while increasing manufacturing costs.\nConversely, a series feed network can reduce efficiency loss and manufacturing difficulty as mentioned above and can solve increase in unit price. However, when the series feed network is used, a transmission signal fed to a radiating element, that is, the phase of electromagnetic waves, also changes according to frequency change, and the direction of a main beam also changes. As a result, gain according to frequency changes dramatically to affect transmission and reception. Especially, in the case of a high gain array antenna, this effect is even larger due to very narrow beam width. Therefore, there is a need for a solution to this problem."}
{"text": "1. Technical Field\nThe present invention relates to a composite membrane and a method of manufacturing the same, and more particularly, to a micropore-containing composite membrane, which can be manufactured to a large area by a simple process, and a method of manufacturing the same.\n2. Description of the Related Art\nMembranes can be classified into porous membranes, dense membranes, and composite membranes, depending upon structures thereof.\nAmong the membranes, composite membranes include an active layer determining separation efficiency, and a porous support capable of minimizing resistance to permeation materials and maintaining pore characteristics and mechanical strength of the composite membranes. The composite membranes may include an active layer of a low thickness to reduce resistance to the fluid passing therethrough, thereby providing selectivity and excellent permeation performance and significantly contributing to practical applications of the membranes due to a merit of adjustability of selectivity and permeability.\nTherefore, studies into manufacturing a novel composite membrane for various applications are being carried out.\nIn the related art, although a method of manufacturing a phase separation-based nano-porous composite membrane is well-known in the art, since it is difficult for a homopolymer blend in a solution or molten state to prevent a phase-separated domain from being gradually coarsened and becoming larger due to high mobility of soft materials, it is difficult to obtain a porous active layer having a pore size of 10 nm or less.\nAs other methods of manufacturing a composite membrane, methods of manufacturing a composite membrane are also known in the art. In these methods, a porous active layer having a pore size of 10 nm or less can be obtained through self-assembly of a block copolymer, chemical vapor deposition of carbon nanotubes, or deposition of nanoparticles. However, in these methods, since it is difficult to manufacture a large-size composite membrane or a composite membrane exhibits low dimensional stability, the composite membrane cannot be used at high temperatures or in various solvents.\nTherefore, the inventors of the present invention have conceived the present invention by paying attention to research results showing that an active layer including pores having a size of 10 nm or less can be manufactured by a simple process using concentration polarization of anorganic sol of reactive, nano-particulate covalent networks on the surface of a porous support and can be applied to a manufacturing process of a large-area composite membrane."}
{"text": "Radio frequency identification (RFID) systems have been proposed for identifying tagged objects for such purposes as taking inventory or tracking movements of objects being transported. Examples are described in commonly assigned U.S. Pat. Nos. 5,300,875; 5,365,551; and 5,448,110.\nRFID systems generally employ a passive or active RF transceiver, called a \"tag\", mounted on each object to be identified or tracked. An interrogator transceiver periodically transmits RF interrogation signals. Upon receiving an interrogation signal, a tag responds by transmitting a response signal containing data which identifies the object and contains any other information which may have been stored or programmed in the tag.\nConventional RFID systems provide little or no interactive feedback in response to actions performed by handling personnel. Specifically, conventional RFID systems lack any means for discriminating in favor of an individual tagged object a human operator is working with at any given moment; instead, conventional RFID systems generally would confuse the operator by providing information regarding all the tagged objects in the vicinity. Furthermore, if a number of personnel are working close to each other, conventional RFID systems cannot direct information about a tag to the specific individual who is handling the tagged object.\nFor example, suppose a number of airport baggage handler personnel are sorting or routing tagged suitcases according to the airline flight destination encoded in a tag attached to each suitcase. Conventional RFID systems lack any means for detecting which individual suitcase a human operator or baggage handler is about to pick up so as to provide to the operator only the destination or routing information for the suitcase that person currently is handling, to the exclusion of information about other nearby suitcases. Presumably because of this and other shortcomings of conventional RFID systems, RFID tags never entered commercial use for tagging airline baggage."}
{"text": "Prior art of this kind are disclosed in Japanese Patent Application Laying-open No. 10-187339 and Japanese Patent Application Laying-open No. 2002-91676.\nIn the prior art of Japanese Patent Application Laying-open No. 10-187339, key patterns are preprinted on the surface of the touch panel, and the touch time during which the user touches the touch panel with a finger is detected. When the touch time is more than a threshold value, the cursor is moved. When the touch time is less than the threshold value, the touch position is detected, and a graphic, number, symbol or the like printed in the position is input.\nAlso, in the prior art of Japanese Patent Application Laying-open No. 2002-91676, when one key (of the character “” (Japanese syllabic character hiragana “tsu”), for example) is pointed at, the graphics related to that (“” (“zu”) and “”, for example) are adjacently displayed. By moving toward these related graphics while pointing at the key and then stopping the pointing, the graphic at the time of stopping the pointing is input.\nAccording to the above mentioned two prior art references, the display position of a graphic, etc. input by touching the touch panel is fixedly decided, and the user cannot freely select the display position.\nTherefore, it is a feature of certain exemplary embodiments to provide a novel information processing apparatus and information input program.\nIt is another feature of certain exemplary embodiments to provide an information processing apparatus and information input program which make it possible to freely select or set a display position of a graphic, etc. input by operating a touch panel.\nIn a first exemplary embodiment, an information processing apparatus that comprises a display, a touch panel provided in relation to the display, a displaying means for displaying a plurality of key patterns and an input area on the display, a touch detecting means that detects a touch-on to or a touch-off from the touch panel at predetermined time intervals, a coordinate position detecting means that detects a coordinate position at a time when the touch detecting means detects the touch-on, a key pattern designation determining means that, when the touch detecting means detects a change of state from touch-off to touch-on, determines whether or not the coordinate position detected by the coordinate position detecting means designates any of the plurality of key patterns, a first coordinate position determining means that, after the key pattern designation determining means determines that some of the key patterns is designated, in response to the detection of a change to touch-off by the touch detecting means, determines whether the coordinate position detected by the coordinate position detecting means immediately before the change to touch-off is a position corresponding to the key pattern or an arbitrary position within the input area, and a first graphic displaying means that, when the first coordinate position determining means determines that the coordinate position at a position corresponding to the key pattern, displays a graphic associated with the key pattern in a predetermined position in the input area, and that, when the first coordinate position determining means determines the coordinate position is an arbitrary position within the input area, displays the graphic associated with the key pattern at the arbitrary position.\nIn the first exemplary embodiment, the information processing apparatus (10: a reference numeral indicative of a corresponding element in the embodiments. The same applies to the following reference numerals.) includes the display (14) and the touch panel (22) provided on the display. A processor, i.e., a CPU core (42) and a step S1 performed thereby function as displaying means. The displaying means (42, S1) causes the display to display an input screen (image) having the plurality of key patterns (66) like a keyboard and the input area (68). The touch detecting means (42, S3, S15, S27) detects at predetermined time intervals that a user touches the touch panel (22) (touch-on) or cancels the touch (touch-off) with a stick (24) or the like. When the touch detecting means detects the touch-on, the coordinate position detecting means (62b, S5, S17) detects a coordinate position at that time in the touch panel. The key pattern designation determining means (S7) determines whether or not the coordinate position at the time of detection of a change of state from touch-off to touch-on by the touch detecting means designates any one of the plurality of key patterns. Then, after the key pattern designation determining means determines that some of the key patterns is designated, when the touch detecting means detects a change to touch-off, the first coordinate position determining means (S29, S39) determines whether the coordinate position detected by the coordinate position detecting means immediately before the change to touch-off is a position corresponding to the key pattern or an arbitrary position within the input area. That is, when some of the key patterns is touched and then touch-off is performed, the first coordinate position determining means determines whether the coordinates at the time of the touch-off is in the position of the key pattern or an arbitrary position not indicative of the key pattern in the input area. Also, when the first coordinate position determining means determines that the coordinate at the time of touch-off is in the position corresponding to the key pattern, the first graphic displaying means (S31, S41) displays a character graphic associated with the key pattern at a predetermined position in the input area (a position subsequent to the previous input position, for example). Conversely, when the coordinate position is an arbitrary position within the input area, the first graphic displaying means displays the character graphic associated with the key pattern at the arbitrary position.\nAccording to the first exemplary embodiment, the coordinate positions at the start (touch-on) and end (touch-off) of an operational input to the touch panel are detected. When the coordinate positions at the time of start and the time of end are on the same key pattern, a character graphic associated with the key pattern is displayed at a predetermined position in the input area. When the coordinate position at the time of start is on the key pattern but it is moved to an arbitrary position within the input area at the time of touch-off, the character graphic associated with the key pattern is displayed at the arbitrary position within the input area. That is, by touching a key pattern through a touch-on operation and moving the stick (24) or the like to an arbitrary position within the input area while keeping it in the touch-on state, it is possible to drag the character graphic for the key pattern into the arbitrary position. Accordingly, by means of this dragging technique, the user can freely select an input position within the input area, i.e., the display position of the character graphic associated with the key pattern.\nA second exemplary embodiment is an information processing apparatus according to the first exemplary embodiment, which further comprises a time measuring means that, when the key pattern designation determining means determines that some of the key patterns is designated, counts a time during which the touch-on state continues, a duration time determining means that determines whether or not duration of the time measured by the time measuring means is equal to or more than a predetermined time, a second coordinate position determining means that, when the duration time determining means determines that the duration time is equal to or more than the predetermined time, determines whether or not the coordinate position detected by the coordinate position detecting means corresponds to the designated key pattern, and a second graphic displaying means that, when the second coordinate position determining means determines that the coordinate position is a position corresponding to the key pattern, displays the graphic associated with the key pattern at a predetermined position in the input area.\nIn the second exemplary embodiment, when the key pattern designation determining means (S7) determines that some of the key patterns is designated at the time of a change to touch-on, the time measuring means (64f, S11) counts the time during which the touch-on state continues. The duration time determining means (S19) determines whether or not the continuing time is equal to or more than a predetermined time. If the duration time is equal to or more than the predetermined time, the second coordinate position determining means (S21) determines whether or not the coordinate position detected by the coordinate position detecting means corresponds to the designated key pattern. If the coordinate position corresponds to the key pattern, the second graphic displaying means (S23) displays the character graphic associated with the key pattern at a predetermined position in the input area.\nAccording to the second exemplary embodiment, when an input of a key pattern is continuously performed for a predetermined time or more, an image associated with the key pattern can be displayed at a predetermined position in the input area. That is, by continuously touching the one and same key pattern, after lapse of a specific time, the character graphic for the same key pattern can be input and displayed at a predetermined position in the input area at each time of touch detection.\nA third exemplary embodiment is an information processing apparatus according to the first exemplary embodiment, which further comprises a valid key setting means that, when the touch detecting means detects a change of touch-off to touch-on, sets the designated key pattern determined by the key pattern designation determining means as a valid key, wherein when the first coordinate position determining means determines that the coordinate position is a position corresponding to the key pattern set as the valid key, the first graphic displaying means displays the graphic associated with the key pattern at a predetermined position in the input area, and when the first coordinate position determining means determines that the coordinate position is an arbitrary position within the input area, the first graphic displaying means displays the graphic associated with the key pattern at the arbitrary position, and the apparatus further comprises a canceling means that, when the touch detecting means detects a change of the touch-on to touch-off, cancels the setting of the valid key.\nIn the third exemplary embodiment, when a change takes place from touch-off to touch-on, the valid key setting means (S9) sets the designated key pattern determined by the key pattern designation determining means as a valid key. Also, when the second coordinate position determining means determines that the coordinate position is a position corresponding to the key pattern set as the valid key, the second graphic displaying means (S31) displays the character graphic associated with the key pattern at a predetermined position in the input area. However, when the second coordinate position determining means determines that the coordinate position is an arbitrary position within the input area, the second graphic displaying means displays the character graphic associated with the key pattern at the arbitrary position. Moreover, when the touch detecting means detects a change from touch-on to touch-off, the canceling means (S35) cancels the setting of the valid key.\nAccording to the third exemplary embodiment, the key pattern indicated at the time of start of an operational input (touch-on) is set as a valid key, and the setting of the valid key is canceled at the end of the operational input (touch-off). Accordingly, even if the input coordinate position is changed by the above mentioned dragging operation, etc. from the key pattern designated at the time of start to another key pattern, it is possible to prevent the other key pattern from being input by mistake.\nThe present invention of a fourth exemplary embodiment is an information processing apparatus according to any one of the first through third exemplary embodiments, which further comprises a predetermined position setting means that, when the first graphic displaying means or the second graphic displaying means displays the graphic associated with the key pattern, sets a position related to the display position of the graphic as the predetermined position.\nIn the fourth exemplary embodiment, when the first graphic displaying means or the second graphic displaying means displays the character graphic associated with the key pattern, the predetermined position setting means (S25, S33) sets the position related to the display position of the character graphic as the predetermined position.\nAccording to the fourth exemplary embodiment, the predetermined position as next character input position is set in the input area.\nA fifth exemplary embodiment is a storage medium storing a program for information input in an information processing apparatus (10) comprising a display (14), a touch panel (22) provided in relation to the display, and a processor (42), the program causing the processor to perform the following steps of: a displaying step (S1) of causing the display to display a plurality of key patterns and an input area; a touch detecting step (S3, S15, S27) of detecting a touch-on to or a touch-off from the touch panel at predetermined time intervals; a coordinate position detecting step (S5, S17) of detecting a coordinate position at a time when the touch-on is detected in the touch detecting step; a key pattern designation determining step (S7) of, when a change of state from touch-off to touch-on is detected in the touch detecting step, determining whether or not the coordinate position detected in the coordinate position detecting step designates any of the plurality of key patterns; a first coordinate position determining step (S29, S39) of, after it is determined in the key pattern designation determining step that some of the key patterns is designated, in response to the detection of a change to touch-off in the touch detecting step, determining whether the coordinate position detected in the coordinate position detecting step immediately before the change to touch-off is a position corresponding to the key pattern or an arbitrary position within the input area; and a first graphic displaying step of, when it is determined in the first coordinate position determining step that the coordinate position is a position corresponding to the key pattern, displaying a graphic associated with the key pattern at a predetermined position in the input area, and, when it is determined in the first coordinate position determining step that the coordinate position is an arbitrary position within the input area, displaying the graphic associated with the key pattern at the arbitrary position.\nIn the fifth exemplary embodiment, the same advantages can be expected as those of claim 1.\nA sixth exemplary embodiment is a storage medium according to the fifth exemplary embodiment, wherein the program further causes the processor to perform the following steps of: a time measuring step (S11) of, when it is determined in the key pattern designation determining step that some of the key patterns is indicated, counting a time during which the touch-on state continues; a duration time determining step (S19) of determining whether or not the duration time counted in the time measuring step is equal to or more than a predetermined time; a second coordinate position determining step (S21) of, when it is determined in the duration time determining step that that the duration time is equal to or more than the predetermined time, determining whether or not the coordinate position detected in the coordinate position detecting step is a position corresponding to the designated key pattern; and a second graphic displaying step (S23) of, when it is determined in the second coordinate position determining step that the coordinate position is a position corresponding to the key pattern, displaying the graphic associated with the key pattern at a predetermined position in the input area.\nIn the sixth exemplary embodiment, the same advantages can be expected as those of the second exemplary embodiment.\nA seventh exemplary embodiment is a storage medium according to the fifth exemplary embodiment, wherein the program further causes the processor to perform a valid key setting step (S9) of, when a change of touch-off to touch-on is detected in the touch detecting step, setting the designated key pattern determined in the key pattern designation determining step as a valid key, wherein when it is determined in the first coordinate position determining step that the coordinate position corresponds to the key pattern set as the valid key, in the first graphic displaying step, the graphic associated with the key pattern is displayed at a predetermined position in the input area, and when it is determined in the first coordinate position determining step that the coordinate position is an arbitrary position within the input area, in the first graphic displaying step, the graphic associated with the key pattern is displayed at the arbitrary position, and the program causes the processor to further perform a canceling step (S35) of, when a change of the touch-on to touch-off is detected in the touch detecting step, canceling the setting of the valid key.\nIn the seventh exemplary embodiment, the same advantages can be expected as those of the third exemplary embodiment.\nAn eighth exemplary embodiment is a storage medium according to any one of the fifth through seventh exemplary embodiments, wherein the program further causes the processor to perform a predetermined position setting step (S25, S33) of, when the character graphic associated with the key pattern is displayed in the first graphic displaying step or the second graphic displaying step, setting the position related to the display position of the character graphic as the predetermined position.\nIn the eighth exemplary embodiment, the same advantages can be expected as those of the fourth exemplary embodiment.\nA ninth exemplary embodiment is an information input method for an information processing apparatus (10) comprising a display (14), a touch panel (22) provided in relation to the display, and a processor (42), which includes the following steps of: a displaying step (S1) of causing the display to display a plurality of key patterns and an input area; a touch detecting step (S3, S15, S27) of detecting a touch-on to or a touch-off from to the touch panel at predetermined time intervals; a coordinate position detecting step (S5, S17) of detecting a coordinate position at a time when the touch-on is detected in the touch detecting step; a key pattern designation determining step (S7) of, when a change of state from touch-off to touch-on is detected in the touch detecting step, determining whether or not the coordinate position detected in the coordinate position detecting step designates any of the plurality of key patterns; a first coordinate position determining step (S29, S39) of, after it is determined in the key pattern designation determining step that some of the key patterns is designated, in response to the detection of a change to touch-off in the touch detecting step, determining whether the coordinate position detected in the coordinate position detecting step immediately before the change to touch-off is a position corresponding to the key pattern or an arbitrary position within the input area; and a first graphic displaying step of, when it is determined in the first coordinate position determining step that the coordinate position is a position corresponding to the key pattern, displaying a graphic associated with the key pattern at the predetermined position in the input area, and, when it is determined in the first coordinate position determining step that the coordinate position is an arbitrary position in the input area, displaying the graphic associated with the key pattern at the arbitrary position.\nIn the ninth exemplary embodiment, the same advantages can be expected as those of the first through the fifth exemplary embodiments.\nAccording to certain exemplary embodiments, by touching a key pattern through a touch-on operation and moving the stick or the like to an arbitrary position in the input area while keeping it in the touch-on state, it is possible to drag the character graphic for the key pattern into the arbitrary position. Accordingly, the user can freely select the input position in the input area, i.e., the display position of the character graphic associated with the key pattern.\nThe above described objects and other objects, features, aspects, and advantages of the exemplary embodiments will become more apparent from the following detailed description of the exemplary embodiments when taken in conjunction with the accompanying drawings."}
{"text": "Computer disk drives are common to personal computers, laptop computers and other computing devices. As with most products, computer disk drives must be shipped from a manufacturer or distributor to another location where the disk drives may be sold or used. As understood by those skilled in the art, computer disk drives are precision electromechanical devices that electronically store data and allow data to be manipulated in accordance with the functioning of the computing devices in which the disk drives are installed. A disk drive includes one or more disks, and at least one read/write component known as the “head” which reads and writes data to and from its corresponding disk. In general, packages or containers for disk drives must provide the proper amount of and support to prevent damage to the drives if the container is dropped or inadvertently contacted by a foreign object. During shipment, a disk drive is shipped with its head in the “landing zone”. The landing zone is an area on the disk where the head is positioned when the disk drive is not in use. Refinement in the construction of some disk drives has resulted in fewer and/or smaller heads being used. Accordingly, for these newer types of disk drives, there is oftentimes not enough friction between the smaller heads and the landing zone to prevent the disk from rotating due to vibration or shock during shipment. Vibrations experienced by these types of disk drives during shipment can cause the disks to partially rotate in repetitive back and forth motions. These small, partial rotations of the disks cause the lubrication to be displaced or separated from the ball bearings and bearing races within the spindle motors that drive the disks. The separation of the lubricant from the bearings and races can result in bearing/race damage due to the lack of proper lubrication. This damage is called “motor fret”. When a disk drive is installed in a computer, the existence of motor fret may be significant enough to cause increased motor noise. Motor noise is a defect that can make the disk drive unsuitable for sale to a consumer.\nThe vast majority of disk drives are shipped from a manufacturer in multi-pack boxes, that is, those boxes/containers that hold at least twenty disk drives. These multi-pack containers are then palletized wherein many containers are packed together and strapped to a pallet.\nExisting multi-pack containers typically include a corrugated outer carton and inner protective inserts that isolate each of the disk drives within the container. Typically, two inserts are used, namely, a top cover and a main insert. The pair of inserts work together as a pair to protect the drive from shock on all axes. One common insert material used is expanded polypropylene, also known as EPP foam. EPP foam is relatively low in cost and durable, as well as resilient enough to provide good shock and vibration protection. The inserts are molded in a desired configuration to hold the set number of disk drives to be packaged within the container. Most multi-pack containers arrange the disk drives on edge and orient them transversely with respect to a long axis of the container. In this arrangement, every disk drive is visible when the upper insert or top cover is removed. Accordingly, each drive can be accessed individually without having to remove other drives. Individual drive access allows bar code scanning, software loading, etc., while minimizing handling of the drives.\nSome structural characteristics common to most if not all inserts include the use of a flat, horizontal plate or floorboard, peripheral vertical side walls that surround the plate or floorboard, and a plurality of partitions or dividers arranged in the space between the vertical sidewalls, gaps or spaces between the dividers forming compartments that receive the disk drives. The purpose of the dividers is to separate each drive from its neighbor so that the drives do not contact one another during shipment. The partitions can be either full or partial height, that is, the partitions can fully cover the drives, or only partially cover the drives thus there being some gap between the inserts. The peripheral edges of the inserts as well as the exposed upper and lower surfaces of the inserts may include a plurality of shock pads or ribs that extend from the floorboard exterior surfaces and contact the inner surfaces of the outer carton that receives the inserts.\nWhen properly designed, the ribs function by compressing to absorb impact energy, and then rebound to essentially their original size and shape. The floorboard itself only serves as secondary cushioning, while primary cushioning is achieved by the externally extending cushioning ribs.\nOne example of a prior art container for multiple disk drives is disclosed in the U.S. Pat. No. 6,588,595. The container of this invention includes three major components, namely, a main insert, a top cover, and a cardboard shell. The main insert is constructed of expanded polypropylene material molded to include a plurality of compartments to receive disk drives loaded therein. The compartments are arranged in one or more rows extending longitudinally along the length of the package. The top cover is placed over the insert and contacts the upper surfaces of the disk drives. A stabilizing member in the form of a central rib communicates with the top cover and extends longitudinally along the row(s) of the disk drives to partially constrain the disk drives during lateral or transverse movement within their respective compartments. This invention has been proven to greatly reduce or eliminate motor fret.\nIn addition to providing shock and vibration protection for the disk drives, it is also desirable to provide a container that may be easily shipped and stored according to international shipping standards. More particularly, one key restriction for palletized loads that may be shipped by air is that the loads must fit into an internationally accepted “cube” size. Typically, the cube is 40 inches×48 inches in length and width, and approximately 45 inches maximum in height. Many multi-pack boxes are designed to fit precisely on pallets to minimize wasted space; however, a need still exists for shipping additional disk drives per pallet in order to further minimize shipping and storage costs. However, effective shock and vibration protection should be maintained if a multi-pack box is modified to accommodate the more cost effective shipping and storage."}
{"text": "Renewable energy sources, such as biofuels, provide a substitute for fossil fuels and a means of reducing dependence on petroleum oil. In light of its low cost and wide availability, biomass is often used as a feedstock to produce pyrolysis oil (which is relatively soluble in water) or bio-oil which, in turn, is used to produce biofuel.\nMany different conversion processes have been developed for converting biomass to bio-oil or pyrolysis oil. Existing biomass conversion processes include, for example, combustion, gasification, slow pyrolysis, fast pyrolysis, liquefaction and enzymatic conversion. Pyrolysis oil is the resultant of thermal non-catalytic treatment of biomass. The thermocatalytic treatment of biomass renders liquid products that spontaneously separate into an aqueous phase and an organic phase. Bio-oil consists of the organic phase. Pyrolysis oil and bio-oil may be processed into transportation fuels as well as into hydrocarbon chemicals and/or specialty chemicals.\nWhile thermolysis processes and other conversion processes produce high yields of such oils, much of the pyrolysis oil and bio-oil produced is of low quality due to the presence of high levels of low molecular weight oxygenates having 5 or less carbon atoms (C5−). Such low MW oxygenates can be in alcohols, aldehydes, ketones, carboxylic acids, glycols, esters, and the like. Those having an isolated carbonyl group include aldehydes and ketones like methyl vinyl ketone and ethyl vinyl ketone.\nSuch oils thus require secondary upgrading in order to be utilized as drop-in oxygen free transportation fuels due to the high amounts of such oxygenates. A known method for converting oxygenates into hydrocarbons is hydrotreating wherein the stream is contacted with hydrogen under pressure and at moderate temperatures, generally less than 750° F., over a fixed bed reactor.\nTransportations fuels predominately contain hydrocarbons having six or more carbon atoms (C6+) (though small amounts of C5 hydrocarbons are present in some gasolines). Thus, hydrocarbons derived by hydrotreating C5− oxygenates are of little value in transportation fuels. Additionally, hydrotreating C5− oxygenates consumes valuable hydrogen in the reactor. Thus, the efficiency of secondary upgrading of pyrolysis oil and bio-oil is compromised by the presence of the C5− oxygenates.\nAlternative processes have therefore been sought for enhancing the efficiency in hydrotreating of oils derived from biomass. Processes for enhancing the yield of hydrotreated pyrolysis oil and bio-oil from streams containing C5− oxygenates, especially C3 to C5 oxygenates, have been sought."}
{"text": "The invention relates to a method of compacting a non woven sintered metal web. The invention also relates to a non woven sintered metal web obtained by said method. Such a non woven metal web may be used in filters and in radiant surface combustion burners.\nBy non woven metal web is meant a web consisting of metal fibers having a diameter of less than 100 micrometer, preferably less than 60 micrometer. e.g. 25 .mu.m, 15 .mu.m, 2 .mu.m or 1.5 .mu.m. These metal fibres have rough outer surfaces and are randomly distributed in the web and are interlocked in an intermingled relationship substantially solely by means of the rough outer surfaces thereof in a frictional engagement. A non woven metal web may comprise two or more layers wherein each layer comprises metal fibres of about the same diameter and wherein the diameters of the different layers are not necessarily equal to each other. By a non woven metal web is also meant a non woven metal web reinforced by one or more wire meshes. A non woven metal web may be produced by a rando-feeder-webber apparatus which is disclosed in GB 1.190.844.\nThis non woven metal web may be sintered or not, but relating to this invention the sintered form is to be preferred. By sintering a bond is created at the points of crossing of the different fibers.\nIt is well known to use non woven metal webs in filters. The pore size is one of the most important characteristics of a filter. It determines the size of the particles which will be kept by the filter. In order to obtain a specified limited pore size the non woven metal web has been compacted by pressing or rolling up to now. However, these conventionally compacted non woven webs have some limitations. The density, i.e. a measure For the compactness (for correct definition, see below), corresponding to a specified pore size is quite high. This causes relatively higher pressure drops over the filter or necessitates the use of larger filter surfaces or the use of ventilators and pumps with more power.\nIt is also known to use non woven metal webs as panels in radiant surface combustion burners. These burners can have a true radiant surface combustion, i.e. no blue flame patches or no free flame, with a thermal output up to 900 kW m.sup.-2. However, until now this has only been possible with non woven webs having a limited surface of about 40.000 mm.sup.2 or less (and this not on a constant basis), or, if one desires to use surfaces greater than 40.000 mm.sup.-2, with thermal outputs which are limited to 500 kW m.sup.-2 or less. This represents a severe restriction for the application of these burners."}
{"text": "The present invention relates to a method and circuit for detecting when the magnitude of a current exceeds a predetermined threshold value.\nCurrent detectors are utilized in myriad of applications; for instance many regulator circuits are seen in the prior art that require the need to provide a regulated output responsive to a detected current exceeding a threshold value. Another application of a current threshold detector may be to control the operating state of high output current operational amplifiers. In present day applications it is a desired goal to produce amplifiers with low quiescent power consumption but such prior art amplifiers have been deficient in output drive current and alternating current (AC) performance. In fact, it is still the goal to provide an operational amplifier that has high drive current capability in order to drive low impedance loads while also having minimum bias drain current to reduce power consumption. In deed, in portable battery powered applications such as mobile telephones, consumer entertainment systems e.g., radios and video games etc., it is especially important to limit the quiescent power consumption of amplifiers used in such applications.\nMost, if not all, present day high output current operational amplifiers are comprised of an input stage that is coupled to an output stage. In response to an alternating input signal being applied to the input stage, the operational amplifier both sources and sinks current to and from a load coupled to the output stage. Typically, the input and output stages are biased at a quiescent drain current to permit quality audio and data processing applications. For example, low power, high output current amplifiers, such as the MC 33178 manufactured by Motorola, Inc., draw approximately 420 microamperes of drain current per amplifier in the quiescent operating mode with no input signal applied thereto. For micro-power and battery powered applications there is a need to reduce drain current requirements over the foregoing while providing amplifiers that can supply high load currents. One manner in which the foregoing may be accomplished is to utilized circuitry for sensing when the input signal is below a certain threshold value to minimize the current drain of the amplifier and once the input signal exceeds this threshold to increase the bias current to the stages so that the amplifier operates in its intended range.\nHence, there is a need for an improved method and circuit that can provide an output signal which may be utilized to control a load by providing a load control signal when a detected current exceeds a predetermined threshold value."}
{"text": "Referring to FIG. 1, a panel 100 used in an electric home appliance is fabricated by bending one side and the other side of a steel sheet.\nThe one side and the other side are arranged adjacent to each other in the bending process and then a corner 100 is formed.\nIn this instance, it is impossible to couple the two sides completely in the bending process and a slit or a hole (hereinafter, a gap region) is usually formed.\nThe gap region 110 happens to expose a sharp surface of the steel sheet and deteriorates safety. Also, the gap region 110 generates tolerance or more errors and it is then difficult to assemble components.\nAccordingly, a welding process is additionally performed to finish the gap region 110 in the prior art to which the present disclosure pertains to. However, such a welding process is likely to raise production cost and the welded portion has to be polished additionally.\nRecently, an insert assembly developed including a wedge 200 and a supporter 300 which are assembled to each other as shown in FIG. 5 and the insert assembly is inserted in the gap region 110.\nHowever, such an insert assembly has to use many components and requires auxiliary assembly lines. The conventional insert assembly has disadvantages of much production time and high production cost."}
{"text": "The use of multiple input multiple output (MIMO) technology has attracted increased attention for use in wireless communications systems because MIMO offers significant increases in data throughput and link range without requiring additional bandwidth or transmit power. The increased performance afforded by MIMO technology stems from higher spectral efficiency (greater number of bits transmitted per second per Hertz of bandwidth), as well as greater link reliability or diversity. Accordingly, MIMO forms an important part of modern wireless communications standards including 3GPP Long Term Evolution (see 3GPP TS 36.213, section Technical Specification Release 10, June 2011, 3rd generationPartnership Project), IEEE 802.11n (WiFi), 802.16 (WiMAX) and HSPA+.\nOne area of concern is the ability to provide robust rank adaptation. Rank adaption refers to the dynamic control of rank according to changing channel conditions. The channel conditions may be determined by such parameters as signal to interference and noise ratio (SINR) and fading correlation between antennae in a MIMO system. With the use of spatial multiplexing, a base station (or eNodeB, or eNB) may send multiple data streams or layers to UEs in a downlink transmission using the same frequency. The number of such layers or streams is defined as the rank. The UE may periodically measure a channel and send a recommendation of the rank to the eNB. The so called rank indicator (RI) may be sent periodically or aperiodically in different schemes. Because the RI reported to the eNB may change with time, the eNB may adjust the number of data streams used in a downlink transmission, based upon the changing RI received from the UE. However, several factors may render this process less than ideal. In some circumstances, the interference levels that may affect channel quality can change substantially between two successive RI reports, in which case, the eNB has no occasion to adjust the rank even though the last reported rank may not be appropriate due to the changed interference conditions. In other circumstances, when a so-called wideband rank is used, the rank indicator reported may be based upon an entire transmission band (wideband), which may be composed of a group of frequency sub-bands used for communications between the UE and eNB. In many cases, the interference conditions may vary substantially between different sub-bands within the wideband, thereby compromising the validity of a wideband RI reported by the eNB for individual sub-bands.\nAnother concern has been raised regarding the use of multiuser MIMO (MU-MIMO) where a UE may transmit precoding matrix indicator/channel quality indicator (PMI/CQI) reports with too high a rank to effectively support the most efficient MU-MIMO scheduling: In MU-MIMO an eNB may schedule multiple different UEs for transmission over the same transmission band. The reporting of an excessively high rank may arise as a consequence of the fact that the UE can only evaluate its own link performance, and in general is not aware of any co-scheduling candidates in the MU-MIMO scheme. As a consequence, even though the link performance to a UE may be maximized by a high rank single user MIMO (SU-MIMO transmission), the system performance may very well be higher if the UE and a second terminal (unbeknownst to the UE) are co-scheduled using lower rank transmissions. The reported information including PMI/CQI/RI may therefore be ill-matched to a MU-MIMO allocation. In particular, the terminal may often select a too high transmission rank to be beneficial for MU-MIMO scheduling.\nIn other circumstances, where the eNB prefers operating under SU-MIMO transmissions, the UE may report a rank of 1, although the eNB may prefer the UE to report an adapted rank, rather than always reporting a rank of 1.\nIt is with respect to these and other considerations that the present improvements have been needed."}
{"text": "1. Field of the Invention\nThe present invention relates to a communication terminal and communication network for partially updating software, a software update method, and a software creation device and method therefor, and more particularly to a software update technique which is capable of downloading and updating a changed part of software embedded in an embedded device.\n2. Description of the Related Art\nEmbedded devices typically refer to devices that software is embedded in chip thereof, such as mobile phones, personal digital assistants (PDAs), etc.\nIn order to update or upgrade (hereinafter simply “update”) software embedded in the embedded device, a user of the embedded device usually visits a business office associated with the embedded device to request update of the software, and an operator of the business office connects the embedded device to an update terminal (e.g., personal computer) having a new version of software to download a new version of the software from the update terminal to the embedded device and change the software to the new version of software.\nIn a conventional method for updating software embedded in the embedded device, the update terminal is first connected with the embedded device to provide a download service for the software embedded in the embedded device. The update terminal then determines whether or not a change has occurred in the software embedded in the embedded device by comparing files of the software with corresponding files of a new version of software stored in the update terminal. Upon determining that a change has occurred in the software embedded in the embedded device, the update terminal updates the software by downloading the new version of software to the embedded device.\nHowever, the above-mentioned conventional software update method has disadvantages in that a lot of software upgrade time is spent. This is because the update terminal has to download all files of the new version of software stored therein to the embedded device."}
{"text": "The present invention relates to non-volatile memory devices and methods for operating them. More particularly it relates to a NAND-type flash memory device capable of shortening a copy-back program time by changing an operation of a page buffer, and a method for operating the page buffer.\nRecently, there has been increasing demand for semiconductor memory devices which do not require periodic refresh operations and are electrically programmable and erasable. A program operation writes data to the memory cells.\nTo achieve high integration of semiconductor memory devices, NAND-type flash memory devices have a plurality of memory cells that share a common connection. In other words, neighboring cells share a drain and source with each other. Unlike NOR-type flash memory cells, NAND-type flash memory cells are capable of reading out information sequentially.\nNAND-type flash memory devices employ a page buffer in order to store large quantities of information or read out stored data within a short time. The page buffer receives a large amount of data from I/O (Input/Output) PAD to provide the data to a memory cell or store the data in the memory cell to output it. In general, page buffers are comprised of a single register so as to temporarily store data. Recently, NAND-type flash memory devices with a dual register have been introduced for the purpose of improving programming speed when programming a large amount of data.\nCopy-back program operation refers to transmitting data stored in a defective cell to a normal cell utilizing page buffers.\nFIG. 1 is a block diagram illustrating a copy-back program operation of a conventional NAND-type flash memory device.\nWith reference to FIG. 1, a conventional copy-back program operation is carried out as follows. Data stored in a defective cell is read out to a first latch unit 24 of a page buffer 20. The data read out from the first latch unit 24 is transmitted to a second latch unit 25. The transmitted data in the second latch unit 25 is programmed to another memory cell that presumably functions properly (or “normal cell”).\nFIG. 2 is a block diagram illustrating program, read, and verification operations of a conventional NAND-type flash memory device.\nWith reference to FIG. 2, if the first latch unit 24 is selected from the first and second latch units 24 and 25, the second latch unit 25 is inactivated, and program operation 51 and read and verification operation 52 are carried out in the first latch unit 24. In contrast, if the second latch unit 25 is selected, the first latch unit 24 is inactivated, and program operation 61 and read, and verification operation 62 are carried out in the second latch unit 25.\nIn the above-mentioned copy-back program operation, there is a high probability that errors occur in transmitting data between the first latch unit 24 and the second latch unit 25. Accordingly, a timing margin may not be sufficiently secured during a copy-back program operation."}
{"text": "Influenza virus is a member of Orthomyxoviridae family (for review, see Murphy and Webster, 1996). There are three subtypes of influenza viruses designated A, B, and C. The influenza virion contains a segmented negative-sense RNA genome. The influenza virion includes the following proteins: hemagglutinin (HA), neuraminidase (NA), matrix (M1), proton ion-channel protein (M2), nucleoprotein (NP), polymerase basic protein 1 (PB1), polymerase basic protein 2 (PB2), polymerase acidic protein (PA), and nonstructural protein 2 (NS2) proteins. The HA, NA, M1, and M2 are membrane associated, whereas NP, PB1, PB2, PA, and NS2 are nucleocapsid associated proteins. The NS1 is the only nonstructural protein not associated with virion particles but specific for influenza-infected cells. The M1 protein is the most abundant protein in influenza particles. The HA and NA proteins are envelope glycoproteins, responsible for virus attachment and penetration of the viral particles into the cell, and the sources of the major immunodominant epitopes for virus neutralization and protective immunity. Both HA and NA proteins are considered the most important components for prophylactic influenza vaccines.\nInfluenza virus infection is initiated by the attachment of the virion surface HA protein to a sialic acid-containing cellular receptor (glycoproteins and glycolipids). The NA protein mediates processing of the sialic acid receptor, and virus penetration into the cell depends on HA-dependent receptor-mediated endocytosis. In the acidic confines of internalized endosomes containing an influenza virion, the HA protein undergoes conformational changes that lead to fusion of viral and host cell membranes followed by virus uncoating and M2-mediated release of M1 proteins from nucleocapsid-associated ribonucleoproteins (RNPs), which migrate into the cell nucleus for viral RNA synthesis. Antibodies to HA molecule can prevent virus infection by neutralizing virus infectivity, whereas antibodies to NA proteins mediate their effect on the early steps of viral replication.\nInactivated influenza A and B virus vaccines are licensed currently as trivalent vaccines for parenteral administration. These trivalent vaccines are produced as monovalent bulk in the allantoic cavity of embryonated chick eggs, purified by rate zonal centrifugation or column chromatography, inactivated with formalin or β-propiolactone, and formulated as a blend of the two strains of type A and the type B strain of influenza viruses in circulation among the human population for a given year. The available commercial influenza vaccines are whole virus (WV) or subvirion (SV; split or purified surface antigen) virus vaccines. The WV vaccine contains intact, inactivated virions. SV vaccines treated with solvents such as tri-n-butyl phosphate (Flu-Shield, Wyeth-Lederle) contain nearly all of the viral structural proteins and some of the viral envelopes. SV vaccines solubilized with Triton X-100 (Fluzone, Sanofi-Aventis; Fluvirin, Novartis) contain aggregates of HA monomers, NA, and NP principally, although residual amounts of other viral structural proteins are present. A live attenuated cold-adapted virus vaccine (FluMist, Medlmmune) was granted marketing approval recently by the FDA for commercial usage as an intranasally delivered vaccine indicated for active immunization and the prevention of disease caused by influenza A and B viruses in healthy children and adolescents, 5-17 years of age and healthy adults 18-49 years of age.\nSeveral recombinant products have been developed as recombinant influenza vaccine candidates. These approaches have focused on the expression, production, and purification of influenza virus type A HA and NA proteins, including expression of these proteins using baculovirus infected insect cells (Crawford et al, 1999; Johansson, 1999; Treanor et al., 1996), viral vectors (Pushko et al., 1997; Berglund et al., 1999), and DNA vaccine constructs (Olsen et al., 1997).\nCrawford et al. (1999) demonstrated that influenza HA expressed in baculovirus infected insect cells is capable of preventing lethal influenza disease caused by avian H5 and H7 influenza subtypes. At the same time, another group demonstrated that baculovirus-expressed influenza HA and NA proteins induce immune responses in animals superior to those induced by a conventional vaccine (Johansson et al., 1999). Immunogenicity and efficacy of baculovirus-expressed hemagglutinin of equine influenza virus was compared to a homologous DNA vaccine candidate (Olsen et al., 1997). Taken together, the data demonstrated that a high degree of protection against influenza virus challenge can be induced with recombinant HA or NA proteins, using various experimental approaches and in different animal models.\nLakey et al. (1996) showed that a baculovirus-derived influenza HA vaccine was well-tolerated and immunogenic in human volunteers in a Phase I dose escalation safety study. However, results from Phase II studies conducted at several clinical sites in human volunteers vaccinated with several doses of influenza vaccines comprised of HA and/or NA proteins indicated that the recombinant subunit protein vaccines did not elicit protective immunity [G. Smith, Protein Sciences; M. Perdue, USDA, Personal Communications]. These results indicated that conformational epitopes displayed on the surface of HA and NA peplomers of infectious virions were important in the elicitation of neutralizing antibodies and protective immunity.\nRegarding the inclusion of other influenza proteins in recombinant influenza vaccine candidates, a number of studies have been carried out, including the experiments involving influenza nucleoprotein, NP, alone or in combination with M1 protein (Ulmer et al., 1993; Ulmer et al., 1998; Zhou et al., 1995; Tsui et al., 1998). These vaccine candidates, which were composed of quasi-invariant inner virion proteins, elicited a broad spectrum immunity that was primarily cellular (both CD4+ and CD8+ memory T cells). These experiments involved the use of the DNA or viral genetic vectors. Relatively large amounts of injected DNA were needed, as results from experiments with lower doses of DNA indicated little or no protection (Chen et al., 1998). Hence, further preclinical and clinical research may be required to evaluate whether such DNA-based approaches involving influenza NP and M1 are safe, effective, and persistent.\nRecently, in an attempt to develop more effective vaccines for influenza, particulate proteins were used as carriers of influenza M2 protein epitopes. The rationale for development of an M2-based vaccine was that in animal studies protective immunity against influenza was elicited by M2 proteins (Slepushkin et al., 1995). Neirynck et al. (1999) used a 23-aa long M2 transmembrane domain as an amino terminal fusion partner with the hepatitis B virus core antigen (HBcAg) to expose the M2 epitope(s) on the surface of HBcAg capsid-like particles. However, in spite of the fact that both full-length M2 protein and M2-HBcAg VLP induced detectable antibodies and protection in mice, it was unlikely that future influenza vaccines would be based exclusively on the M2 protein, as the M2 protein was present at low copy number per virion, was weakly antigenic, was unable to elicit antibodies that bound free influenza virions, and was unable to block virus attachment to cell receptors (i.e. virus neutralization).\nSince previous research has shown that the surface influenza glycoproteins, HA and NA, are the primary targets for elicitation of protective immunity against influenza virus and that M1 provides a conserved target for cellular immunity to influenza, a new vaccine candidate may include these viral antigens as a protein macromolecular particle, such as virus-like particles (VLPs). Further, the particle with these influenza antigens may display conformational epitopes that elicit neutralizing antibodies to multiple strains of influenza viruses.\nSeveral studies have demonstrated that recombinant influenza proteins could self-assemble into VLPs in cell culture using mammalian expression plasmids or baculovirus vectors (Gomez-Puertas et al., 1999; Neumann et al., 2000; Latham and Galarza, 2001). Gomez-Puertas et al. (1999) demonstrated that efficient formation of influenza VLP depends on the expression levels of viral proteins. Neumann et al. (2000) established a mammalian expression plasmid-based system for generating infectious influenza virus-like particles entirely from cloned cDNAs. Latham and Galarza (2001) reported the formation of influenza VLPs in insect cells infected with recombinant baculovirus co-expressing HA, NA, M1, and M2 genes. These studies demonstrated that influenza virion proteins may self-assemble upon co-expression in eukaryotic cells."}
{"text": "Automotive experts quite typically recommend that the squeegees of windshield wiper assemblies be replaced at least annually as a consequence of their deterioration due to abrasion, aging, etc. To minimize the expense of replacement of an entire windshield wiper system including arm, pressure applying superstructure, and squeegee assembly, only the latter is replaced. In the usual case, the squeegee assembly includes both the squeegee itself and a flexor or backing which is applied to the existing pressure applying superstructure mounted on the wiper arm.\nOf course, some means must be provided whereby the squeegee assemblies may be easily disassembled from and replaced on the superstructure which also provides for positive retention of the squeegee assembly on the superstructure during use.\nAs a consequence, there have been a large variety of proposals for retainer clips to accomplish this purpose. Representative of such attempts, and constituting the most pertinent prior art known to the applicant are the following: U.S. Pat. No. 3,626,544 issued Dec. 14, 1971 to Lopez; U.S. Pat. No. 3,879,794 issued Apr. 29, 1975 to Roberts; U.S. Pat. No. 3,919,736 issued Nov. 18, 1975 to Bourassa; U.S. Pat No. 4,009,503 issued Mar. 1, 1977 to Sharp; U.S. Pat. No. 4,156,951 issued June 5, 1979, also to Sharp; U.S. Pat. No. 4,388,742 issued June 21, 1983 to Kimber; and U.S. Pat. No. 4,389,746 issued June 28, 1983 to Riester. Also of possible relevance are French patent publication No. 2,443,354 and German Gebrauchsmuster No. E82 20 739.9. Of the foregoing, U.S. Pat. No. 4,388,742 to Kimber may be the most relevant.\nThe retainers of each of the foregoing patent publications have a variety of advantages and accompanying disadvantages which will be readily appreciated by those skilled in the art. In general, it is desirable that the clip be easily handled, both during an assembly and a disassembly process, that it has sufficient strength as to not be broken or permanently distorted during installation or removal of a squeegee assembly, that it have universal applicability in terms of accommodating superstructures of various configurations and the ability to mount squeegee assemblies of varying length, that it provide for positive retention and firm securement of the squeegee assembly to the superstructure and that it provide a high degree of retention in the proper placement during the assembly process.\nIn general, the prior art suffers from the deficiency of being unable to fulfill one or more of the foregoing desirable attributes of a retaining clip. The present invention is directed to providing a retaining clip that fulfills all of the foregoing desirable attributes."}
{"text": "1. Field of the Invention\nThe present invention relates to an electric rotating machine used as a motor or an alternator in a vehicle.\n2. Description of Related Art\nGenerally, a vehicle-mounted electric rotating machine includes a rotor rotatable around a rotating shaft, a stator core formed with slots arranged along the circumferential direction of the stator core so as to face the outer periphery of the rotor, and a stator coil constituted of a plurality of coil wires wound around the slots. The coil wire includes in-slot portions accommodated in the slots which are different in their circumferential positions, and turn portions each connecting two adjacent ones of the coil wires. The turn portions form a coil end at each of both axial ends of the stator coil, the coil end projecting axially outward from the end surface of the stator coil.\nSince the stator coil of the electric rotating machine having such a structure increases in temperature when it operates, various countermeasures for the temperature increase have been proposed. For example, Japanese Patent Application Laid-open No. 2001-195302 (Patent document 1) describes installing an oil jacket on each of both end surfaces of the stator core so as to surround the coil end, providing an oil chamber between the oil jacket and the end surface, and passing coolant to the oil chamber to cool the coil end.\nFurther, Japanese Patent Application Laid-open No. 2004-180376 (Patent document 2) describes disposing, above the coil end, a gutter including a coolant supply opening through which coolant to be supplied to the surface of the coil end is taken and a guide connected to the coolant supply opening to guide coolant discharged from the coolant supply opening in order that coolant is guided to a desired direction to efficiently cool the coil end.\nIn the case of Patent document 1, the coolant passage is a closed space, and when the oil chamber is filled with coolant, the entire surface of the coil end becomes wet by coolant so that it is cooled. However, the cooling structure described in Patent document 1 needs a coolant supply pump, because the coolant passage has a relatively large fluid resistance. This increases manufacturing cost.\nIn the case of Patent document 2, since coolant flowing in from the coolant supply opening and flowing out from the guide flows down on the surface of the coil end by the action of the gravity, it is likely that some parts of the surface of the coil end do not become wet, causing the cooling to be non-uniform. Generally, the envelope surface of a coil end of a stator of an electric rotating machine is not flat, and accordingly, the non-uniform cooling tends to occur easily."}
{"text": "Lifting and transportation of loads from one place to another is a regular routine in various industries including textile, mechanical, civil, cargo, automotive, the plastic industry, the woodworking industry, the foodstuffs industry and the drinks industry. Load lifting ventures involves transporting loads manually or with the help of load lifting devices. The task of carrying loads of various shapes, sizes and weights manually by laborers is a physically challenging and time consuming process.\nThere are various assemblies available in the prior art for lifting purposes like an assembly having clamping jaws connected by a rod and moving on a column, the lower jaw carries the load and the upper jaw is connected through a lever to an operating lever which is reciprocated to cause the load to climb the column, the jaws gripping the column alternately. Springs urge the jaws to the clamping position. The load may be raised either by inserting a lever between the jaws or for heavier loads a lever is pivoted to the lower jaw and contacts a roller on the upper jaw. Another assembly concerns a linear sliding guide comprising a shaped rail which extends in the longitudinal direction having sliding surfaces, and a slider which is axially displaceable on the sliding surfaces of the shaped rail and has sliding bearings which bear against the sliding surfaces.\nThe assemblies available in the prior art have a fixed load carrying capacity and their overall structure utilizes a lot of workspace. Another drawback is that these assemblies are difficult to assemble, mount and require a lot of maintenance. Also, they lack the flexibility in terms of lifting the loads of various sizes and are expensive to manufacture. Further, they exhibit a lot of friction which hampers their performance.\nThus, there is a need for a new assembly for lifting and transporting the loads which are easy to assemble, to mount, is flexible, inexpensive to manufacture and that can be used in a variety of operations depending upon the type of loads being lifted."}
{"text": "Computer-driven route planning applications are utilized to aid users in locating points of interest, such as particular buildings, addresses, and the like. Additionally, in several existent commercial applications users can vary a zoom level, thereby enabling variation of context and detail as a zoom level of a map is altered. For example, as a user zooms in on a particular location, details such as names of local roads, identification and location of police and fire stations, identification and location of public services, such as libraries, museums, and the like can be provided to the user. When zooming out, the user can glean information from the map such as location of the point of interest within a city, state, and/or country, proximity of the point of interest to major freeways, proximity of the point of interest to a specific city, and the like. In some applications, satellite images can be utilized to provide users with additional detail regarding a particular geographic location or region. For example, a prospective purchaser of a house can obtain an overhead satellite image of the house, thereby enabling the prospective purchaser to view lines of occupation, proximity of the house to other adjacent houses, and other information that can be pertinent to the user.\nFurthermore, conventional computer-implemented mapping applications often include route-planning applications that can be utilized to provide users with directions between different locations. Pursuant to an example, a user can provide a route planning application with a beginning point of travel and an end point of travel (e.g., beginning and ending addresses). The route planning application can include or utilize representations of roads and intersections and one or more algorithms to output a suggested route of travel. These algorithms can output routes depending upon user-selected parameters. For instance, a commercial route planning application can include a check box that enables a user to specify that she wishes to avoid highways. Similarly, a user can inform the route planning application that she wishes to travel on a shortest route or a route that takes a least amount of time (as determined by underlying algorithms). Over the last several years, individuals have grown to rely increasingly on route planning applications to aid them in everything from locating a friend's house to planning cross-country road trips."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a water conservation apparatus and, more particularly, to a master control system for conserving water by sprinkler systems within a geographical region.\n2. Description of the Prior Art\nFresh water is an increasingly valuable commodity. Indeed, many regions in short supply of fresh water are required to purchase and import water much as any other required commodity. Municipalities, as well as entire counties and states, have taken and continue to take drastic steps in order to conserve water. Regulations to conserve water have been enacted to affect industrial, commercial and residential users. Examples include enacting prohibitions against the use of water for non-essential purposes such as filling of swimming pools. Fairly common ordinances by many municipalities or other governmental subdivisions regulate the use of residential sprinkler systems in some fashion. Commonly, residents can only use their sprinkler systems on certain days (e.g., odd or even). Large commercial users of water, such as car washes, have been required to recycle the water that they use. Such regulatory schemes are helpful but not totally effective. Some residential users simply ignore the regulations. Others conceal their use by sprinkling their lawns, for example, at night when such use is more difficult to detect. Also, such regulatory schemes tend to implement simple rules (e.g., odd-even) that do not provide the flexibility or versatility to provide a greater degree of control by which the numbers of users dispensing water at any given time can be easily controlled.\nIn U.S. Pat. No. 4,209,131, a computer-controlled irrigation system is disclosed for automatically irrigating an agricultural area having a plurality of crop-growing fields. The system includes a central station and a number of remote stations located in the fields. A computer is programmed to define a predetermined schedule for delivery of pre-established quantities of water to each field. The delivery of water can be made dependent on environmental conditions such as wind, speed, temperature, humidity or other factors. The system interrogates the remote stations for flow meter, water pressure meter, environment conditions and alarm information, to diagnose proper operation and to monitor the delivery of water. The system, therefore, is arranged to meter appropriate amounts of water to irrigate fields to ensure that the crops are provided with appropriate amounts of water.\nIn U.S. Pat. No. 5,813,606, a radio controlled sprinker system is disclosed for use with a residential sprinkler system that is provided with a transmitter unit that includes at least one moisture sensor. The moisture sensor monitors the moisture in the ground and generates a signal representative of the moisture level in the ground. Activation signals are transmitted on the receipt of a moisture signal which indicates that the moisture level is below a desired amount in the region of the transmitter. Therefore, the local transmitters effectively enable the sprinkler systems only when the ground conditions call for additional water.\nIn U.S. Pat. No. 4,838,310, a remotely controlled irrigation system is disclosed used in conjunction with golf course, business parks and along state highways to irrigate plants in various remote areas that are difficult to manually control. The patent is primarily concerned with providing power at the remote locations for operating a receiver that can be actuated by a central transmitter. This is achieved by providing a turbine or impeller in the path of the pipe feeding the sprinkler system which rotates when fluid moves through the pipe to operate a generator which charges a local battery.\nIn U.S. Pat. No. 4,626,984, a remote computer control for irrigation systems is disclosed for controlling a number of different irrigation systems. Each individual system has its local controller for controlling an individual system. The patent is primarily concerned with the central computer which would normally shut down the entire system. The patent discloses the use of intelligent remote units that can process and execute stored programs. By distributing intelligence among the remote units, the system becomes less dependent on the operation of the central computer. Should the central computer fail, each intelligent remote station will continue to operate its associate irrigation system in accordance with instructions stored therein.\nAn irrigation control system is disclosed in U.S. Pat. No. 4,396,149, which has as one of its objectives the conservation of water and energy by providing optional or highly efficient irrigation to growing plants. Soil moisture sensors transmit data which is used to regulate the operation of the irrigation system and to provide water on a need basis.\nTherefore, while the technologies have been available, municipalities and other local and regional water authorities have not utilized such technologies in a combination with a plurality of individual residential sprinkler systems for selectively disabling such sprinkler systems. While the prior art has been primarily concerned with enabling water dispensing systems on a need basis, such systems have not been provided the means for overriding independently operated sprinkler systems in order to disable them and selectively prevent water from being discharged through such systems."}
{"text": "1. Field of the Invention\nThe present invention relates to a transfer apparatus which displays an image recorded in digital form by a digital still camera (DSC), a video camera, a personal computer or the like through a transmission type image display device formed by a liquid crystal display device (LCD), and transfers the displayed image to a photosensitive recording medium such as an instant photographic film which develops color by light, thereby forming an image.\n2. Description of the Related Art\nConventionally known examples of a method for transferring (i.e., printing) or recording a digitally-recorded image to or on a photosensitive recording medium include an ink jet system using a dot-type printing head, a laser recording system, and a thermal recording system.\nA printing system like the ink jet system has various problems. For example, printing takes time, ink is likely to cause clogging, and precision printing results in the sheet being moistened by ink. The laser recording system involves an expensive optical component such as a lens, resulting in high apparatus cost. Further, the laser recording system and the thermal recording system require considerable power consumption, and are not suited to be carried about.\nThus, generally speaking, the transfer apparatuses used in these systems and, in particular, the ones used in the ink jet system have a problem in that the more precise the apparatus, the more complicated the driving mechanism and the control mechanism, and the larger and the more expensive the apparatus, printing taking a lot of time.\nIn this regard, JP 10-309829 A and JP 11-242298 A disclose transfer apparatuses of the type in which a display image is formed on a photosensitive recording medium like an instant film by using a liquid crystal device, thereby achieving simplification in structure and a reduction cost.\nThe electronic printer disclosed in JP 10-309829 A is capable of copying the display screen of a liquid crystal display on a photosensitive medium to produce a hard copy of a quality equal to that of a photograph. However, in order to copy the display screen of the liquid crystal display on the photosensitive medium in this electronic printer, an optical component such as a rod lens array is arranged between the display screen of the liquid crystal display and the photosensitive medium, so that a predetermined distance (total conjugate length) is required between them. In the example shown, the requisite distance is 15.1 mm. Further, the optical component is rather expensive.\nIn the case of the transfer apparatus disclosed in JP 11-242298 A, there is no need to use an expensive optical component such as a lens or to secure an appropriate focal length. Thus, as compared with the conventional transfer apparatuses, a further reduction can be achieved in terms of size, weight, power consumption, and cost. As shown in FIG. 7, a photosensitive film 400 is closely attached to the display surface of a transmission type liquid crystal display (hereinafter referred to as LCD) 300, and a light source (back light 100) provided on the opposite side of the photosensitive film 400 with respect to the LCD 300 is turned on. That is, a fluorescent lamp 101 is switched on to turn on the back light, whereby the image displayed on the LCD 300 is transferred to the photosensitive film 400.\nFurther, as shown in FIG. 8, the above-mentioned publication discloses another embodiment, according to which a lattice 200 is provided between the back light 100 and the LCD 300, whereby diffusion of light from the back light 100 is restrained. That is, the light is approximated to parallel rays. Further, by providing a spacer 201 consisting of a rectangular hollow member between the lattice 200 and the LCD 300, it is possible to prevent the image of the frame of the lattice 200 (the shadow due to the frame) from being taken by the photosensitive film 400, thus improving the clarity of the image formed on the photosensitive film 400 to a satisfactory degree from the practical point of view without providing an optical component or securing an appropriate focal length.\nFurther, as shown in FIG. 7, the publication discloses an example of a transfer apparatus in which the thickness of the LCD 300, that is, the sum total of the thicknesses of the following components: a polarizing plate 301 on the display surface side, a glass substrate 302, a liquid crystal layer 303, a glass substrate 304, and a polarizing plate 305 on the back light 100 side is 2.8 mm and in which the image on the screen of the LCD 300 with a dot size of 0.5 mm is transferred to the photosensitive film 400. To prevent diffusion of the light from the LCD 300, there is provided a 5 mm lattice with a thickness of 10 mm, and a 20 mm spacer 201 is arranged between the lattice 200 and the LCD 300. Further, the LCD 300 and the photosensitive film 400 are closely attached together to effect image transfer without involving blurring (unclarity) of the image.\nIn this case, an image displayed with a dot size of 0.5 mm is transferred with an enlarge dot size of up to 0.67 mm, which means an enlargement by approximately 0.09 mm on one side, and yet the image obtained is satisfactory from the practical point of view.\nAs described above, in the transfer apparatus disclosed in JP 11-242298 A, image transfer is effected, with the liquid crystal display (LCD) and the photosensitive film being closely attached together, to prevent blurring (unclarity) of the image and to obtain an image satisfactory from the practical point of view. It is to be noted, however, that exposure of the photosensitive film in this arrangement involves the following problems.\nFirst, as shown in FIG. 7, on the outermost surface of the LCD 300, there is arranged the film-like polarizing plate 301, which is closely attached to the photosensitive film 400 during exposure. When the photosensitive film 400 is moved to perform a post-processing, the photosensitive film 400 and the polarizing plate 301 are rubbed against each other to thereby flaw the film-like polarizing plate 301, and the flaw on the polarizing plate 301 is transferred to the photosensitive film 400. Further, this flaw causes scattering of light, resulting in deterioration in the image quality.\nIt might be possible for the polarizing plate and the photosensitive film to be closely attached together during exposure and slightly spaced apart from each other when the photosensitive film is moved. For this purpose, however, it would be necessary to provide, apart from the photosensitive film moving mechanism, a mechanism for effecting close attachment and detachment of the photosensitive film, which is contradictory to the requirement for a reduction in cost and size.\nFurther, generally speaking, a photosensitive film, for example, an instant film, which is the easiest to use, is kept in a lightproof case until it is loaded in a transfer apparatus. Since this lightproof case is equipped with an opening frame somewhat larger than the film, the following procedures must be followed before the photosensitive film can be brought into close contact with the polarizing plate.\nFirst, prior to exposure, one photosensitive film is extracted singly from the lightproof case, and brought into close contact with the surface of the polarizing plate on the surface of the LCD. In this condition, exposure is performed, and, after the completion of the exposure, the photosensitive film is separated from the polarizing plate surface, and moved for a next processing (In the case of an instant film, a processing liquid tube provided in the film sheet is pushed open).\nThese procedures must be repeated for each photosensitive film. In particular, separating the photosensitive film from the polarizing plate surface does not square with automation (or mechanization).\nRecently, the screens of LCDs have progressed in terms of definition, and LCDs with an increased number of pixels and a smaller dot size are being commercialized. For example, as LCDs using low-temperature polysilicon type TFTs, UXGA (10.4 inches; 1200xc3x971600 pixels), XGA (6.3 and 4 inches; 1024xc3x97768 pixels) are on the market.\nAn attempt to apply an LCD with such a high-definition screen to the transfer apparatus disclosed in JP 11-242298 A would lead to the following problem. In the case of UXGA, the dot size of each of the RGB pixels is approximately 0.04 mm on the shorter side. In a transfer apparatus as disclosed in the above-mentioned publication, in which enlargement in dot size is involved, it would be impossible to transfer an LCD image of such a minute dot size to a photosensitive film with satisfactory clarity in a condition in which the dots of the RGB pixels are clearly distinguishable.\nIt is an object of the present invention to eliminate the above problems in the prior art and to provide a transfer apparatus which can realize a substantial reduction in size, weight, power consumption, and cost with a simple structure and which can also be formed as a portable device.\nAnother object of the present invention is to provide a transfer apparatus which is applicable to various types of liquid crystal displays ranging from a liquid crystal display of an ordinary pixel density to a liquid crystal display with a high definition screen having a high pixel density and which makes it possible to obtain a photographic image of a desired degree of clarity, from a photographic image which is satisfactory from the practical point of view to a high-definition photographic image of a higher level of clarity.\nTo achieve the above objects, the present inventors have conducted careful study on a transfer apparatus which makes it possible to obtain a photographic image of a desired degree of clarity, which is of higher practical value, and which allows use of a transmission type image display device, such as a liquid crystal display, which has a high-definition screen of a high pixel density in a structure in which the liquid crystal layer is held between two sets of substrates and polarizing plates. As a result of the study, the present inventors have found that, to prevent blurring (unclarity) of the image, which is inevitably generated when bringing the transmission type image display device and the photosensitive recording medium out of contact with each other, that is, when separating them from each other to achieve a higher practical value with a simple structure, it is necessary to set the sum total of the thicknesses of the substrate and the polarizing plate on the photosensitive recording medium side of the transmission type image display device in accordance with the separation distance between the two components.\nThe present invention provides a transfer apparatus comprising a light source, a transmission type image display device in which a liquid crystal layer is held between two sets of substrates and polarizing plates and a photosensitive recording medium wherein the light source, the transmission type image display device and the photosensitive recording medium are arranged in series along a direction in which light from the light source advances, and a display image transmitted from the transmission type image display device is transferred to the photosensitive recording medium, and wherein the transmission type image display device and the photosensitive recording medium are arranged in a non-contact state, and a distance between the transmission type image display device and the photosensitive recording medium and a sum total of thicknesses of a substrate and a polarizing plate at least on a side of the photosensitive recording medium in the transmission type image display device are set in accordance with a definition of the display image.\nPreferably, the sum total is not more than 1.0 mm.\nPreferably, the distance is 0.01 mm to 3 mm.\nPreferably, the display image and the image transferred to the photosensitive recording medium are substantially identical in size.\nPreferably, each pixel size of the image display device is not more than 0.2 mm.\nIt is preferable that the transfer apparatus further comprises a substantially parallel rays generating element arranged between the light source and the image display device.\nPreferably, the substantially parallel rays generating element comprises a porous plate having a plurality of through-holes, and wherein the porous plate has a thickness not less than three times the diameter or equivalent diameter of the plurality of through-holes.\nPreferably, the plurality of through-holes are parallel to each other and have a circular or polygonal cross section."}
{"text": "The present invention relates to a method of treating Poly(p-phenylene-2,6-benzobisoxazole) (PBO) fiber to reduce the effects of heat and moisture on the fiber.\nPoly(p-phenylene-2,6-benzobisoxazole) (PBO) fiber is a high strength and high modulus fiber. However, PBO fiber degrades rapidly when exposed to moisture and heat. Methods of modifying the surface of PBO fibers have previously been developed. For example, Chinese Patent Document No. CN 101235590 discloses dipping the PBO fiber into a chemical solution and then ultrasonically treating the fiber to modify the surface of the fiber. However, there is a further need for a method of reducing the detrimental effects of moisture and heat on PBO fibers.\nElimination or reduction of the detrimental effects on PBO fiber caused by exposure to moisture and heat would provide a stable high performance fiber that may be used in lightweight armor systems. Improvements in the mechanical properties of PBO fibers through secondary chemical process modifications of base level fibers, would allow for a treated and improved PBO fiber that can be incorporated into lightweight soft armor systems and Fiber Reinforced Plastics (FRPs) for combat helmets and backings to ceramic for small arms protection."}
{"text": "This invention relates to a control system for a step motor, and more particularly to a closed-loop control for controlling a step motor in response to rotation of the step motor.\nA closed-loop control for a step motor is already known wherein a rotational angular position of a step motor is detected and a field winding of a next phase is energized at an optimum point of time, that is, at an optimum rotational angular position of the step motor to drive the step motor in order to prevent step-out of the step motor. In determining such an optimum point of time, also a rotational speed of the step motor is commonly taken into consideration.\nWhere a load is relatively heavy or is to be driven at a high speed, the closed-loop control technique is preferably employed. One of application of the closed-loop control is, for example, a daisy wheel typewriter or printer in which a rotary type wheel commonly called a daisy wheel is rotated to position a type of a selected character to a print position to allow printing of the character. In the typewriter or printer of this type, the rotor of the step motor is rotated from one to another rotational angular position in one or the other direction in response to a character selecting instruction. As the time required for rotation from one to another position reduces, apparently the overall printing speed of the typewriter or printer will rise. Accordingly, the closed-loop control is commonly used in high speed typewriters or printers.\nA concept of an angle of lead may be involved in order to attain a high speed rotation of a step motor under a closed-loop control. It is known that where the lead angle of a step motor is between 1 step and 2 steps, that is, electrically between 90 and 180 degrees, the step motor will rotate continuously like a brushless motor. In this instance, theoretically an angle of lead of 1.5 steps or electrically 135 degrees will yield a maximum rotational speed whereas the rotational speed will be in the minimum when the lead angle is 1 step or 2 steps, that is, electrically 90 or 180 degrees, and the speed will vary continuously between 1 step and 2 steps, that is, between 90 and 180 degrees.\nConventionally, a rotary encoder such as an optical encoder is used to detect a rotational position and speed of a step motor. Step motor controls including a rotary encoder are normally very expensive compared with controls without a rotary encoder because such a rotary encoder is expensive.\nMeans is also known which detects an electromotive force (EMF) or back EMF induced in a field winding of a step motor to detect a rotational angular position of the rotor of the motor. This technique is disclosed, for example, in U.S. Pat. Nos. 4,282,471 and 4,480,218. In a step motor control disclosed in U.S. Pat. No. 4,282,471, a back EMF which is induced in a non-excited phase field winding by a mutual inductance with an excited phase field winding is detected to detect a rotational position of the rotor of the motor. On the other hand, in a step motor control disclosed in the former patent, EMF which is induced in a field winding by a permanent magnet of the rotor of a step motor is detected directly by a specifically provided sensing means to detect a rotational position of the rotor.\nFeedback signals obtained from a rotary encoder or from back EMF are used in most cases to determine whether or not a step motor is actually at a proper rotational angular position or whether or not a step motor is in a step-out condition so that, when the step motor is in a step-out condition, either a positional error may be compensated for or the motor may be restored to its home position. Otherwise, feedback signals are used to determine when a next phase field winding is to be energized. In the latter case, a timing of next energization may be selectively determined in response to the feedback signal and clock pulse signals of the system. Accordingly, infinite velocity control cannot be attained."}
{"text": "The use of metallic super-elastic alloys, such as Ni—Ti (nitinol) and other bi- or tri-metal alloys, has been documented in a variety of technical applications, including fasteners, connectors, gaskets, clamps and seals. Many such uses have required temperature in order to activate the material and change its physical state, while others have used mechanical forces that impart stress to cause a super-elastic physical deformation in the material. Of particular concern to the instant inventor is the application of the super-elastic material to connectors. The use of non-corrosive, metallic super-elastic material offers a decided advantage in high performance connecting assemblies, versus more conventional connectors requiring threaded fasteners, springs, clamps or other holding or securing mechanisms. Particularly it can withstand more wear than alloys used in conventional connectors due to its harder surface characteristics. It can also withstand extreme vibrations and not loosen due its elastic preloaded condition without using conventional adhesives to hold the assembled components and/or the connector itself together. Adhesives used with conventional connectors make them very difficult to disassemble, whereas it is generally possible to make a super-elastic connector completely reversible.\nU.S. Pat. Nos. 5,395,193 and 5,584,631 to Krumme et al., discuss the use of nickel-titanium shape memory retainers in an optimized elastic condition that have super-elastic or pseudo-elastic properties. These fasteners are said to be useful for eyeglass assembly; they are placed onto a pin to retain components together.\nIt is commonly known that nitinol (formally known as Nickel Titanium Naval Ordinance Laboratory, but other super-elastic nickel-titanium alloys being included in this definition) tubing, wire or rod can be used as a mechanical drive shaft. The use of metallic super-elastic alloys, such as Ni—Ti (nitinol) and other bi- or tri-metal alloys, has been documented in a variety of technical applications, including fasteners, connectors, gaskets, clamps and seals. Many such uses have required temperature in order to activate the material and change its physical state, while others have used mechanical forces that impart stress to cause a super-elastic physical deformation in the material. Of particular concern to the instant inventor is the application of the super-elastic material to connectors. The use of non-corrosive, metallic super-elastic material offers a decided advantage in high performance connecting assemblies, versus more conventional connectors requiring threaded fasteners, springs, clamps or other holding or securing mechanisms. Particularly it can withstand more wear than alloys used in conventional connectors due to its harder surface characteristics. It can also withstand extreme vibrations and not loosen due its elastic pre-loaded condition without using conventional adhesives to hold the assembled components and/or the connector itself together. Adhesives used with conventional connectors make them very difficult to disassemble, whereas it is generally possible to make a super-elastic connector completely reversible.\nU.S. Pat. No. 5,683,404 to Johnson, entitled “Clamp and Method for its Use”, further discusses shape memory materials that are “pseudo-elastic”, defining these materials to be super-elastic, because of their ability to exhibit super-elastic/pseudo-elastic recovery characteristics at room temperature. In other words, a material is super-elastic when, if sufficient stresses are applied, such materials exhibit martensitic activation (i.e., deform from an austenitic crystal structure to a stress-induced structure postulated to be martensitic in nature), returning thence to the austenitic state when the stress is removed. The alternate crystal structures described give the alloy super-elastic or pseudo-elastic properties. Poisson's Ratio for nitinol is about 0.3, but this ratio significantly increases up to approximately 0.5 or more when the shape memory alloy is stretched beyond its initial elastic limit it is at this point that stress-induced martensite is said to occur. i.e., the point beyond which the material is permanently deformed and thus incapable of returning to its initial austenitic shape. A special tool is employed by Johnson to impart an external stretching force that deforms the material which force is then released to cause the material to return to its original condition. While the device is stretched, a member is captured by it and securely clamped when the stretching force is released. This device is intended for use in clamping and does not contemplate traditional connecting operations of the kind addressed by the present invention. Another use envisioned by Johnson is in connecting the modular components of a medical device, as described in his U.S. Pat. No. 5,858,020, by subjecting a thimble component made of shape memory material to an external stretching stimulus to elongate and thereby reduce its transverse dimension. Upon release of the stretching force, this component returns towards its original rest dimension, contacting and imparting a force on another component. This is a sequential stretching and relaxation of the super-elastic material rather than a simultaneous activation and retention operation. Also, special structures are necessary on the thimble to allow the stretching force to be imparted.\nIn U.S. Pat. No. 5,197,720 to Renz, et al., a work piece is held within a clamping tool by an expansion element made of shape memory material that is activated by mechanical force. In this way, torque is transmitted through the shape memory member. This device is useful for bringing parts together for holding the work piece in order to perform an operation. It does not, however contemplate a use as a connector. U.S. Pat. No. 5,190,546 to Jervis discloses insertion into a broken bone cavity of a split member made of shape memory material using a SUPER-elastic alloy. The split member holds the walls of the bone cavity when radial compressive forces acting on it are released. In order for the radial compressive force to reduce the diameter, the component must be split, allowing the reduction in dimension for insertion. It does not act as an interposed member in a connecting assembly.\nOthers have sought to utilize the properties of shape memory materials as locking, connector and bearing elements, e.g. U.S. Pat. No. 5,507,826 to Besselink et al, U.S. Pat. No. 5,779,281 to Kapgan, et al., and U.S. Pat. No. 5,067,827 to Arnold, respectively; however, such approaches have required temperature to be applied during use. U.S. Pat. No. 5,277,435 to Kramer, et al. and U.S. Pat. No. 5,876,434 to Flomenblit, et al. similarly has relied upon temperature to activate the shape memory effect. Such dependence on extrinsic activation by temperature introduces an added process step and may further be disadvantageous in certain other applications.\nU.S. Pat. No. 5.842.312 to Krumme, et al., entitled, “Hysteretic Damping Apparati and Methods”, employs shape memory tension elements to provide energy dissipation. Such elements can be placed between building structures, etc. which are subject to vibration, serving to absorb the energy created by their relative movement. However, this patent does not contemplate the vibration dampening effect of a SUPER-elastic material in the formation of a connector.\nNitinol is especially useful for transmitting torque while in a bowed or bent shape. These types of drive shafts have proven useful in orthopedic surgical applications where drilling or reaming of curved bones is necessary. One application is to use a drill or reamer with a nitinol drive shaft to clean out the center of a femoral bone before implanting a prosthesis or femoral nail. These bones typically have a bow with a 90-inch radius and require a flexible reamer for the procedure. Nitinol tubing can be used for this application since it is cannulated and can be passed over a guide wire that is placed down the femur before the reaming process begins. Since the tubing is solid it is very easy to clean after the surgical operation since there are no crevices for blood to get trapped in. Earlier designs utilized spring drive shafts and cleaning was extremely difficult since blood could get trapped between the windings of the spring. The earlier spring designs also had difficulties when run in the reverse direction since springs tend to be strong while being used in one direction, however when run in the opposite direction they tend to unwind. To prevent this unwinding problem several manufacturers have added an additional spring inside of the primary spring, which is wound in the opposite direction. Since one spring is inside of the other this contributes to the difficulties with cleaning and further obviates the need for an alternative shaft design.\nWith the market demand increasing for these novel nitinol drive shafts there have been many attempts to develop safe coupling methods for attachment to the shaft. One difficulty that engineers have been faced with is presented when nitinol tubing exceeds its torsional or fatigue stress limits; it has been known to fail catastrophically and fragment into several sharp pieces. This is dangerous when inside of a patient and poses severe concerns if these types of products are to be used reliably. Historically there have been no solutions offered to limit the stress in the drive shaft, which would eliminate the presently lingering concerns over breakage during use.\nAnother difficulty is presented with the attachment of the fittings to the nitinol drive shaft. The connection must be reliable and not create any unnecessary stress on the tubing. This will lead to early failure of the shaft. Typically in the orthopedic reamer example mentioned above one end of the nitinol tube has a stainless steel fitting which attaches to a power instrument and on the opposite end either a stainless steel reamer head or an intermediate modular fitting that connects to a reamer head. Several attempts to create reliable attachments have been made.\nOne approach to the above-mentioned attachment problem has been to use an epoxy to “glue” the fittings onto the nitinol shaft. However, temperatures in the sterilization process and the criticality of surface preparation have rendered this approach unreliable.\nAnother approach has been to attach the fitting to the shaft with a laser weld; however, the welding process embrittled the tubing and it was known to fail torsional demands in testing. A cross hole and pin were placed through the fittings, however this added approach further proved useless since the matching hole in the tubing created a tremendous stress riser in the tubing causing failure at very low torsional values. In the example mentioned it was known to fail anywhere between 2 to 4 N–M.\nYet another approach has been to press fit the nitinol shaft into a fitting with approximately 0.002-inch interference. Initial trials worked, however when put through rigorous fatigue tests the tubing placed too much hoop stress on the fittings causing them to fail rather than the shaft. The solution proposed to fix that problem was to add a long section on the fitting that was loosely fit around the tube. This would allow the stress to transition slowly into the area where the press fit was done. This worked successfully, however the solution created a need for the fitting to be extremely long in comparison to the reamer heads being used. This is undesirable since the reamer must follow the curvature of the bone and it did not solve the issue of limiting the torque in the shaft to ensure the safety of the drive shaft during use.\nThus, there remains a primary need to provide a coupling system that is safe and effective for use in surgical and industrial applications where flexible drive shafts are necessary.\nAccordingly, there is a need to form a connecting assembly using a durable metallic, non-corrosive connector assembly, which are simple to install using relative motion to activate the assembly.\nThere is a further need to form a secure connection between components that minimizes the micro-motional wear characteristics of the assembly, enhancing its useful life.\nThere is another need to form a fastened assembly that does not require temperature for its activation.\nThere is still a need to form an assembly using a fastener that adjusts for differences in thermal coefficients of expansion or contraction of dissimilar materials comprising those components being fastened.\nThere is still a further need for a connector with elastic properties that allow more forgiving tolerances during manufacturing of the assembly components.\nThere is another need to use a nitinol drive shaft to replace the spring drive shafts in orthopedic instruments and many industrial tools to simplify the cleaning process and ensure consistent torque resistance in the forward and reverse directions.\nThere is a further need to provide a coupling system which will limit the torque in the nitinol drive shaft fitting to ensure that the coupling limits the torque before the tubing breaks.\nThere is yet a need for a coupling that will not place any unnecessary stress on the nitinol tubing causing it to prematurely break.\nThere is also a need to shorten the length of the fitting so that the reamer can follow the natural contour of the inside of the bone while transferring the stress smoothly so as to ensure the strength of the fitting."}
{"text": "1. Field\nThe following description relates to an optical imaging system including lenses.\n2. Description of Related Art\nA plurality of optical imaging systems are mounted in a portable terminal, such as a camera, a mobile device, or a tablet. For example, the optical imaging systems are mounted on each of a front surface and a rear surface of the portable terminal.\nThe optical imaging system mounted on the rear surface of the portable terminal is used to photograph a subject positioned at a far distance from the portable terminal. Conversely, the optical imaging system mounted on the front surface of the portable terminal is used to photograph a subject positioned at a near distance from the portable terminal. In particular, the optical imaging system mounted on the front surface of the portable terminal generates an optical distortion phenomenon such as a barrel distortion. Therefore, there is a need to develop an optical imaging system capable of reducing the optical distortion phenomenon appropriate to photograph a subject at a short distance."}
{"text": "The present invention relates to radio frequency (RF) power amplifiers (PA) module. Portable devices such as laptop personal computers (PC), Personal Digital Assistant (PDA) and cellular phones with wireless communication capability are being developed in ever decreasing size for convenience of use. Correspondingly, the electrical components thereof must also decrease in size while still providing effective radio transmission performance. However, the substantially high transmission power associated with RF communication increases the difficulty of miniaturization of the transmission components.\nA major component of the wireless communication device is the radio frequency PA. The PA is conventionally in the form of a semiconductor integrated circuit (IC) chip or die in which signal amplification is effected with substantial power. The amplifier chip is interconnected in a circuit with certain off-chip components such as inductors, capacitors, resistors, and transmission lines used for controlling operation of the amplifier chip and providing impedance matching of the input and output RF signals. The amplifier chip and associated components are typically assembled by interconnected metal circuit and bond wires on a printed circuit board (PCB) having a dielectric substrate or a lead frame.\nAmong significant considerations in the miniaturization of RF amplifier circuits is the required impedance matching for the input and output RF signals of the amplifier. Input and output impedance matching circuits typically include capacitors, resistors, and inductors in associated transmission lines or micro strips for the RF signals into and out of the amplifier chip. However, these impedance matching circuits may require specifically tailored off-chip components located remotely from the amplifier IC chip. Accordingly, the application circuitry must include many electrical input and output terminals or bonding Pins to which the corresponding portions of the off-chip impedance matching circuits are separately joined. This increases the difficulty of assembly and required size of the associated components, and affects the overall manufacturability of the portable devices.\nOne important requirement for the state-of-the-art wireless devices is to provide power amplification in a plurality of frequency bands. The quality and power level of the amplified RF signals need to be properly controlled. The amplification of RF signals is required to be linear over a wide signal power range in each of the plurality of frequency bands. Preferably the amplification is reduced or increased according to input RF signal, transmittance range and data rate so that power consumption can be optimized."}
{"text": "Minimally invasive surgery is a surgical approach that involves use of instruments inserted through several tiny incision openings to perform a surgery causing minimal tissue trauma.\nThis minimally invasive surgery relatively reduces changes in metabolism of the patient in the period of post-surgical care, so it is beneficial to rapid recovery of the patient. Therefore, using such minimally invasive surgery shortens length of a hospital stay of the patient after the surgery and allows patients to return to normal physical activities more quickly. In addition, minimally invasive surgery causes less pain and reduces scar to patients after surgery.\nThe most general form of the minimally invasive surgery is endoscopy. Among them, a laparoscopy that involves minimally-invasive inspection and operation inside abdominal cavity is known as the most general form of endoscopy. To operate the standard laparoscopic surgery, an abdomen of the patient is insufflated with gas, and small incisions (about ½ inch or less) are formed for use as an entrance of a tool for the laparoscopic surgery, through which a trocar is inserted. In general, laparoscopic surgical tools include a laparoscope (for observation of a surgical site) and other working tools. Here, the working tools are similar in structure to the conventional tools used for small incision surgery, except that the end effector or working end of each tool is separated from its handle by an elongated shaft. For instance, working tools may include a clamp, a grasper, scissors, a stapler, needle holder, and so forth. To perform the surgery, a user, such as a surgeon, puts the working tool into a surgical site through the trocar, and manipulates it from the outside of abdominal cavity. Then, the surgeon monitors the procedure of the surgery through a monitor that displays the image of the surgical site that is taken by the laparoscope. The endoscopic approach similar to this is broadly used in retroperitoneoscopy, pelviscopy, arthroscopy, cisternoscopy, sinuscopy, hysteroscopy, nephroscopy, cystoscopy, urethroscopy, pyeloscopy, and so on.\nThe modes for these minimally-invasive surgeries may be classified into the single port surgery mode and the multi-port surgery mode based on the number of ports to insert surgical instruments into the surgical field, etc.\nWhen a surgery is performed in the multi-port surgery mode, the surgery may rather be easy despite of disadvantageous increasing number of incisions, on the other hand, when the surgery is performed in the single port surgery mode, the level of difficulty of the surgery may rise due to the problems such as collisions among the surgical instruments whereas the reduced number of incisions makes the surgery less invasive. In this context, the applicant of the invention has already disclosed minimally-invasive surgical instruments with multiple degrees of freedom in the Korean patent application No. 2008-51248 and No. 2008-61894, and the applicant has further disclosed the advantages of the single port surgery mode and minimally-invasive surgical instruments suitable for the same in the Korean patent application No. 2008-79126 and No. 2008-90560, the contents of which are incorporated herein by reference in its entirety.\nTherefore, in case of real surgery, it is preferable to allow selectively implementing both of the single port surgery mode and the multi-port surgery mode according to various surgical situations in a single surgical robot system."}
{"text": "As well known in the art, exemplary nuclear reactors in South Korea include System Plus 80 from Combustion Engineering (CE)(e.g. Hanbit Nuclear Power Plant (NPP) Units 3 and 4), Korean Standard Nuclear Reactors (e.g. Hanul NPP Units 3, 4, 5, and 6; Hanbit NPP Units 5 and 6), OPR1000 (e.g. new Kori NPP Units 1 and 2), APR1400 (new Kori Units 3 and 4, new Ulgin Units 1 and 2, UAE NPP Units 1 to 4), 2-loop Pressurized Water Reactor, etc.\nSuch NPPs (referred hereinafter simply to as “Nuclear Reactors”) include reactor coolant systems (RCSs) in reactor containment buildings as also disclosed in Korean Patent No. 10-1473665 entitled “Tubing Support Apparatus for Replacement of Parts of Nuclear Power Plant”.\nThe RCS includes a reactor containing atomic piles and at least one heat-transfer circuit connected thereto.\nThe circuit includes a steam generator and at least one coolant pump that circulates coolant between the steam generator and the reactor.\nIn addition, the circuit includes a pressurizer that allows the temperature and pressure of the coolant to be kept constant.\nA first large-diameter pipe or a hot leg is connected to one side of the reactor and one side of a suction part of a coolant chamber in the steam generator so as to transmit the coolant, which is heated with contact with a core of the reactor, to the steam generator.\nFurther, a circulation pipe called a cross-over leg connects one side of an outlet of the coolant chamber and one side of an inlet of a swirl chamber in the coolant chamber.\nA cold leg connects the swirl chamber in the coolant chamber and the reactor. The coolant that is cooled at the steam generator and drawn by the coolant pump is transmitted to the reactor via the circulation pipe and the cold leg to cool the core.\nIn such a nuclear reactor, the steam generator 1 has a typical substructure as shown in FIG. 1.\nThat is, the steam generator 1 of the nuclear reactor 1 includes a stay cylinder 10 that is maintained at high temperature and a cylindrical skirt 20 that supports the stay cylinder from a lower section thereof.\nThe skirt 20 is fixedly supported at the bottom thereof by a sliding base 30 by means of a plurality of stud bolts 40. The sliding base 30 is supported by a plurality of (e.g. four) semi-spherical sliders 52, which is provided on a forged bolted plate 50 so as to accommodate a slight motion occurring during the operation of a nuclear reactor.\nIn this structure, the stud bolt 40 has a conventional solid bolt structure.\nDuring operation, the steam generator produces high temperature heat that, when transmitted to the lower side sliding base 30, elevates the temperature of the sliding base 30 as shown in a heat analysis thermal distribution diagram of FIG. 2a. \nThe heat analysis thermal distribution diagram shows that, as a result of analysis with respect of an insulation state, an operation temperature, material, heat transfer, an air flow around a steam generator, etc, a dead air region 70 defined by the skirt 20 and the sliding base 30 is heated to high temperature of up to 131° C. through convection and radiation of high temperature heat (300° C. or more) from the steam generator 1. It could be seen that in such a high temperature state, the skirt 20 was subjected to thermal deformation (e.g. thermal expansion of up to 2.4 mm).\nThe thermal deformation of the sliding base 30 also causes serious problems as follows.\nThe thermal deformation, such as thermal expansion, of the sliding base 30 causes restriction to a free motion of the sliding base 30 or interference with an upper surrounding structure of the steam generator 1, resulting in structural vibration of the steam generator 1 and the coolant pump.\nThe structural vibration causes pipe wearing and vibration stress to the steam generator 1 and after a long operation time, causes material fatigue of small-diameter pipes that have borate embrittlement due to accumulated borate, leading to leakage of a boric acid solution.\nFurther, as a construction factor of the plumbing in a nuclear reactor, due to weld contraction occurring by a final connection welding between an intermediate pipe and a steam generator nozzle, residual stress remains in the sliding base and a vertical support for a pump. Due to the residual stress, subsidence of the sliding base 30 occurs in response to the weld contraction. As a result, after installation of the steam generator and after final connection welding, the sliding base generally sinks by about 1 mm or less.\nSuch subsidence of the sliding base may be considered as residual stress of the sliding base 30, and it increases friction force with respect to the sliding base at an initial operating stage of a nuclear reactor, interfering with a transverse sliding motion of the sliding base.\nFurther, thermal deformation of the sliding base also causes the operating steam generator to be tilted, which may lead to misalignment of parts even after cold shutdown of the steam generator.\nConsequently, if the sliding base is deformed so as to be inclined, the steam generator is accordingly tilted so that deformation and interference occur to support structures for the coolant pump and the steam generator, thereby further increasing vibration stress of a nuclear reactor.\nIt is reported from many countries that such vibration stress causes wear of pipes of the steam generator and of internal components of the coolant pump, and fatigue failure of tubing connected to the RCS during operation of nuclear power plants.\nIf a nuclear reactor is operated for a long period of time in such condition, a leak may occur from a mechanical seal of a coolant pump and small-diameter pipes in the nuclear reactor.\nAccordingly, there is a need to develop a technique to prevent thermal deformation of the sliding base 30 of a nuclear reactor and resultant wear of small pipes of the steam generator and coolant system in a nuclear reactor occurring due to structural vibration of the nuclear reactor.\nThe applicant has proposed three solutions to address this problem.\nA first solution is an air-circulation sleeve 80 that is installed through the center of a sliding base 30 supporting a steam generator to naturally circulate air using a venture effect. This enables a dead air region 70 defined by a stay cylinder and a skirt 20 of the steam generator 1 to be cooled to effectively prevent high temperature heat from being transferred to the sliding base 30 from the steam generator 1, thereby preventing vibrations of a nuclear reactor due to thermal expansion of the sliding base 30.\nA second solution is an air-circulating shim plate (not shown) between the bottom of the skirt 20 supporting the steam generator 1 and an engaging surface of the sliding base to allow ambient air to be introduced into a dead air region 70 defined by a stay cylinder and a skirt 20 of the steam generator 1 to be cooled to effectively prevent high temperature heat from being transferred to the sliding base 30 from the steam generator 1, thereby preventing vibrations of a nuclear reactor due to thermal expansion of the sliding base 30.\nA third solution is a heat insulation support plate 90 that is closely attached to a lower portion of the steam generator near the stay cylinder and the skirt. The heat insulation support plate is composed of a heat insulation material, a heat shield panel, and a plurality of rigid pieces. The heat insulation support plate is thus prevented from sagging due to operating vibration of a nuclear reactor. The heat insulation support plate serves to effectively prevent high temperature heat from being transferred to the sliding base 30 through the dead air region 70, thereby preventing vibrations of a nuclear reactor due to thermal expansion of the sliding base 30.\nThe above solutions proposed by the applicant contributes to suppression of thermal deformation of the sliding base 30, having a great effect of stable operation of a nuclear reactor and improvement in lifecycle of equipment of the nuclear reactor.\nHowever, the above solutions are techniques that are applicable at the time of replacement of a steam generator after one cycle of a commercial operation of a nuclear reactor.\nThus, the solutions are difficult to be applied to normally operated nuclear reactor. Further, since even in the operated nuclear reactor, damage of equipment such as a heat pipe or the like can be prevented only when causes of thermal deformation are previously removed, there is a great need to develop a technique for cooling a substructure of a steam generator in a nuclear reactor, wherein the technique is applicable even to a nuclear reactor that is under construction, at a test run stage, or just before replacement of a steam generator."}
{"text": "As the material for a permanent magnet formed of an oxide, there is known hexagonal system M type (magnetoplumbite type) Sr ferrite or Ba ferrite. The ferrite magnetic materials made of the ferrites are made available as permanent magnets in the form of ferrite sintered bodies or bonded magnets. In recent years, with miniaturization and enhancement of performances of electronic components, permanent magnets made of ferrite magnetic materials have also been required to have high magnetic characteristics while being small-sized.\nAs the indexes of the magnetic characteristics of a permanent magnet, generally, the residual magnetic flux density (Br) and the coercive force (HcJ) are used. Those which are high in these are evaluated as having high magnetic characteristics. Conventionally, from the viewpoint of improving the Br and the HcJ of the permanent magnet, a study has been done by the changing of the composition such as the inclusion of a prescribed element in a ferrite magnetic material.\nFor example, in the following Patent Document 1, there is shown oxide magnetic materials capable of providing a ferrite sintered magnet having high Br and HcJ by allowing M type Ca ferrite to include at least La, Ba, and Co.\nFurther, in the following Patent Document 2, there are disclosed oxide magnetic materials capable of providing ferrite sintered magnets having high Br and HcJ by allowing M type Ca ferrite to include La, Sr, and Co. Further, in the following Patent Document 3, there are disclosed sintered magnets having high Br and HcJ by allowing M type Sr ferrite to include Sr, La, and Co."}
{"text": "1. Field of the Invention\nThe present invention relates to an electrophotographic photoreceptor. The invention also is concerned with a method of producing the electrophotographic photoreceptor.\n2. Description of the Prior Art\nVarious organic and inorganic electrophotographic photosensitive materials have been proposed, among which a dual-layer type photoreceptor having a charge generating layer and a charge transport layer is becoming a matter of concern as a photoreceptor having high photosensitivity. In this type of photoreceptor, the charge generating layer is made of a charge generating compound such as a phthalocyanine pigment-type organic compound as disclosed in Japanese Patent Unexamined Publication No. 59-15253 or inorganic compounds of tellurium-arsenic-selenium alloy type as disclosed in Japanese Patent Unexamined Publication No. 50-15137, as well as various other known compounds. On the other hand, various compounds are usable as the charge-transport substance, such as compounds of of poly-N-vinylcarbazole as shown in Japanese Patent Unexamined Publication No. 52-77730 or of pyrazoline derivative type as disclosed in Japanese Patent Unexamined Publication No. 49-105537. These charge generating and charge transport compounds have been actually put into practical use in photoreceptors.\nMeanwhile, a high-speed electrophotographic printers making use of a laser beam have been put into practical use. This type of printer is capable of performing printing at a high speed, while reducing the size of the light source unit and, hence, the size of the whole printer, particularly when a laser diode is used as the light source. It is also to be noted that organic photoreceptors are disposable by users themselves.\nThe known organic photoreceptors, on the other hand, has weak durability to abrasion as compared with inorganic photosensitive materials In consequence, the photoconductive layer is worn down to reduce the thickness thereof, so that the electrophotographic characteristics cannot have a long life. In order to overcome this problem, it has been proposed to form a protective film on the surface of the photoreceptor to protect it from mechanical damaging force For instance, Japanese Patent Unexamined Publication No. 58-83857 discloses a photoreceptor in which a selenium photosensitive layer is covered by a protective layer made of an electron donor and a binder. In this photoreceptor, a polycarbonate resin, polyester resin or a polyurethane resin is used as the binder used in the protective layer. Japanese Patent Unexamined Publication Nos. 61-22345 and 61-27550 propose photoreceptors having protective layers made of alkylether melamine formaldehyde. The protective layers of this material, however, tend to increase the level of the residual potential.\nIn use, a corona discharge makes the surface of the photosensitive material to cause the surface thereof to hold electric charges. Therefore, the photoreceptor, particularly when it is of organic type, tends to exhibit a degradation (oxidation) at the surface thereof during long use, increasing a tendency for the surface to be wetted with moisture. The wetted surface of the organic photosensitive material tends to allow an easy leakage of electric charges, resulting in a lowered resolution of printed images and other defects.\nPhotoreceptors using fluorine compounds are also known. For instance, Japanese Patent Unexamined Publication No. 62-206559 discloses an electrophotographic photoreceptor having a surface layer containing a fluorine oligomer compound. Japanese Patent Unexamined Publication No. 55-7762 discloses a photoreceptor made of a material containing a fluorine-containing compound. Japanese Patent Unexamined Publication No. 61-95358 discloses an image forming member provided with a resin layer containing fluorine polymer. Japanese Patent Unexamined Publication No. 58-23031 discloses an image carrier containing a surfactant having long-chain fluorinated alkyl group. An image carrier having a surface layer made of a material including a fluorine-containing silane coupling agent is proposed in Japanese Patent Unexamined Publication No. 61-205950.\nThe known organic electrophotographic photoreceptors mentioned above generally exhibit inferior wear-resistance characteristics due to small surface hardness as compared with inorganic photosensitive materials. In order to overcome this problem, Japanese Patent Unexamined Publication Nos. 58-83857, 61-22345 and 61-27550 propose to use resins having high levels of surface hardness. Such a countermeasure, however, is still unsatisfactory.\nAnother problem encountered with the known organic photoreceptors is that surfaces of these materials tend to be wetted due to absorption of water as a result of deterioration (oxidation) of the surfaces caused by a corona discharge during the use, with the result that the resolution of the print image is lowered.\nKnown photoreceptors containing fluoride component also are still unsatisfactory.\nArts shown in Japanese Patent unexamined Publication Nos. 62-206559 and 61-95358 cannot provide required strength of the coating layer because the mechanical strength of the organic coating layer is lowered as a result of use of a fluorine compound as the material of the binder of the organic coating material.\nArts shown in Japanese Patent Unexamined Publication Nos. 55-7762, 58-23031 and 61-205950 cannot provide sufficient lubrication effect because the fluorine compounds used therein are composed of perfluoroalkyl groups having short fluorine chains."}
{"text": "Technical Field\nEmbodiments disclosed herein relate to a method of and loom for teaching Braille by weaving with a loom, wherein the levers to move the harnesses correspond to the keys on a Perkins Brailler (a Braille typewriter).\nThe National Federation for the Blind estimates that 1.1 million people are blind in the United States. Fortunately the blind have several options these days for reading and recording written material. In the past, only learning Braille, a system of six raised dots invented by Louis Braille, enabled one to read words, and mathematics or music symbols. FIG. 1 illustrates the alphabet in Braille. FIG. 2 illustrates the numbers 0–9 in Braille. Words do not always have to be Brailled letter for letter. Many words have contracted or short-form representations in Braille. An excellent reference of commonly used Braille representations of words is The Braille Enthusiast's Dictionary by Alan J. Koenig and M. Cay Holbrook, which is incorporated herein by reference.\nDavid Abraham, a wood-working teacher at the Perkins School for the Blind invented a Braille typewriter in 1951, called the Perkins Brailler. As illustrated in FIG. 3, the Perkins Brailler is configured with six keys and a spacer bar: three keys to the left of the spacer bar and three keys to the right of the spacer bar. FIG. 3 is a top view of a Perkins Brailler. The three keys to the left of the spacer bar make the raised dots in positions three, two, and one, (from left to right), respectively. The three keys to the right of the spacer bar make the raised dots in positions four, five, and six (from left to right), respectively. Each of the six keys is just under ¾ inch wide, and the spacer bar is 1¾ inches wide. The keys are 1 inch apart between centerlines, and the centerlines of keys #1 and #4 are 1½ inches from the centerline of the spacer bar.\nRecently technology advances have enabled electronic word files to be transposed into a spoken voice, which “reads” the text file to the user. A person may also create an electronic data file by use of available software that transcribes the spoken word of the user into a text file. With such advances, the incentive to learn Braille has diminished, and as a result, so has the Braille literacy rate.\nTwenty years ago in the 1980's forty percent of the blind could read Braille. Today, Braille literacy rates have dropped to ten percent. Because many Braille resources still exist and Braille is in public use in, for example, ATMs, or elevators, the blind still need to know Braille. Additionally, the blind need the ability to quickly create their own Braille documents, and thus the skill of typing in Braille is desirable to have. There is a need for a fun way to learn how to type Braille, which also allows the student to practice his or her knowledge of Braille.\nThe other background for understanding this invention is weaving on looms. Weaving is accomplished by passing a crosswise thread in a predetermined perpendicular pattern under and over lengthwise threads. Looms increase the efficiency of weaving by providing a loom frame that supports a plurality of lengthwise threads (called warp threads). The warp threads are threaded through loops (called heddles) made of wire or string, which may be attached between top and bottom rods or bar members of a harness. When the harness is raised or lowered, it accordingly raises or lowers the attached heddles the same distance, which pull the warp threads threaded through the heddles apart from the other warp threads, creating a sometimes triangular-cross-sectioned space (called a shed) through which to pass the crosswise thread (called the weft). This eliminates the need to move the weft up and over and down an under warp threads. Instead, the weft thread is wound on a shuttle which is passed from one hand to the other through the shed created by the separated warp threads. A good reference book of weaving, which is incorporated by reference, is Ward and Weave by Robert LeClerc available from LeClerc Looms of Canada.\nAt least two harnesses are needed to weave: one to raise a set of warp threads, the other to raise the remaining set of warp threads. However, most looms have at least four harnesses. Commonly, the number of harnesses available on looms increases in multiples of four.\nEach harness is attached to an actuator. It is the actuator that the loom operator moves in order to raise or lower the attached harnesses, thereby separating the selected set of warp threads from the remaining warp threads. In weaving vernacular, if the actuator is mounted up top or on the side, it is referred to as a “lever.” If it is down below the warp, it is referred to as a treadle. A thin, flexible and strong member, such as cording, is often used to tie an actuator to one or more harnesses. Looms that may be used on a table are referred to as table-top looms and most often have top or side levers. Floor looms, so named because they are set up on the floor, most often have treadles that are connected to the harnesses and the treadles are normally operated by foot rather than hand.\nLeClerc Looms makes a portable, table-top loom called the Voyager. It is sold with or four and eight harnesses. Each harness is referred to as a “shaft,” and thus the Voyager is offered in both a four and eight “shaft” model. An eight “shaft” model is illustrated in FIG. 4. A loom frame 20 may include a top board 22 with eight levers 24–38. Each lever is rotatable about a bar from the back of top board 22 to the front of top board 22 and is securable by magnetic forces operating between the magnets 40a and 40b: one magnet 40a in the lever and one magnet 40b directly below it in the top board 10. Each lever is tied to a harness 42–56. Each harness has heddles through which the warp threads 58 are threaded. By rotating a lever, the harness 42–56 to which it is attached is raised, pulling the desired warp threads apart from the remainder. A shuttle 60 may then be passed from one side to the other through a shed (not shown) formed by the separated warp threads (not shown). LeClerc and other loom making companies, such as, for example, Schacht, Louet, and Ashford, all sell four- or eight-“shaft” (or even twelve), table-top looms with levers up top or on the side."}
{"text": "The present invention relates to a new and distinct cultivar of Calibrachoa plant, botanically known as Calibrachoa sp., hereinafter referred to by the cultivar name ‘Dancaldarpink’.\nThe new Calibrachoa cultivar is a product of a planned breeding program conducted by the Inventor, Gabriel Danziger, in Moshav Mishmar Hashiva, Israel. The objective of the breeding program was to develop a new Calibrachoa cultivar having attractive flower colors, a desirable plant habit, and a strong, vigorous growth habit.\nThe new Calibrachoa cultivar originated from a cross made in a controlled breeding program by the Inventor in 2003 in Moshav Mishmar Hashiva, Israel. The female or seed parent is a Calibrachoa sp. designated as ‘CV.217’ (unpatented). The male or pollen parent is a Calibrachoa sp. designated as ‘CV.100’ (unpatented). ‘Dancaldarpink’ was discovered and selected by the Inventor in February of 2003 as a flowering plant within the progeny of the stated cross in a controlled environment in Moshav Mishmar Hashiva, Israel.\nAsexual reproduction of the new cultivar by stem tip cutting was first performed in May of 2003 in Moshav Mishmar Hashiva, Israel, and has demonstrated that the combination of characteristics as herein disclosed for the new cultivar are firmly fixed and retained through successive generations of asexual reproduction. The new cultivar reproduces true-to-type."}
{"text": "1. Field of the Invention\nThe invention relates to a device for adhesive attachment of a cover, configured for receiving a binding strip, to the flanks of a spine of a book block which is formed of bound printed sheets and transported by means of a circulating transport device, wherein the cover to be combined with a binding strip is transferred, in accordance with the cycle of the transported book block and approximately in the same direction to the book block\"\"s transport direction, to the spine of the book block for positioning thereat and is connected with the book block during the subsequent pass through a pressing path with pressing device.\n2. Description of the Related Art\nIt is known to provide, for the purpose of reinforcing, the spine of a book with strips of paper, textile materials or the like, a so-called binding strip, wherein the binding strips are separated from a web, which extends transversely or longitudinally to the path of back gluing processing, and are then combined with a cover being fed to the book block. Parallel glue strips are applied onto the cover during the course of feeding at a greater spacing to one another than the thickness of the book block in order to connect by means of a conveying roll, arranged downstream of the cover feeding device, the cover to the binding strip and the binding strip to the book block, i.e., to its spine provided with a glue layer, while the book block and the cover are being moved in the same direction.\nIt would be possible also to provide the underside of the binding strip with two spaced-apart glue strips.\nThe further transport of the cover connected with the binding strip is carried out by the transport device of the back gluing machine for transporting the book block where subsequently a pressing of the cover and of the binding strip onto both flanks, provided with glue, of the book block spine is carried out. For this purpose, the book block spine, which has been provided at the adjoining flanks with strips of glue, passes through a pressing station in which the cover and the laterally projecting parts of the binding strip are engaged by a parallel-moving pressing device and are pressed against the book block (disclosed in xe2x80x9cTechnologie der Klebebindungxe2x80x9d by Alfred Furler, particularly in paragraphs 8.5 and 8.6, but also well-known from the back gluing machine of the Muller Martini Bookbinding Systems).\nFrom the time of feeding the cover up to the point of pressing the cover and the binding strip onto the book block, there is no certainty that the cover can be reliably connected via the glue strips to the binding strip, in particular, because the parts of the binding strip projecting laterally past the spine of the book block and the oppositely positioned cover have no planar surface and no means have been provided so far to secure them in their relative position to one another. This results in the condition that the cover and the binding strip on the transfer path into the pressing station with the pressing device can be moved relative to the book block.\nIt is an object of the present invention to provide a method and a device for the adhesive attachment of a cover provided with a binding strip onto an book block by which method and device the reliability of a precise connection of the cover with the book block is ensured.\nIn accordance with the present invention, this is achieved in that at least one of the lateral parts of the binding strip projecting laterally past the book block spine are loaded on the side facing away from the cover by a pressing element.\nIn this connection, a mechanical device, for example, a holding-down device, can provide a beneficial and simple pressing means.\nOne embodiment of a device for performing the method according to the invention comprises a cover conveying device and a binding strip conveying device positioned alongside a transport device for the book block and having arranged downstream thereof a pressing device acting on the cover, wherein the device is characterized in that the cover conveying device and the binding strip conveying device have correlated therewith a pressing element in the form of a compressed air supply device or an adjustable holding-down device acting on the lateral parts of the binding strip projecting laterally past the book block.\nPreferably, the pressing element acting on the lateral projecting parts of the binding strip can be realized by compressed air.\nExpediently, the compressed air can be supplied in a controlled fashion, i.e., can be switched on and off and/or can be adjusted with respect to amount and speed, i.e., can be adjusted to the respective conditions.\nFor reinforcing the adhesion between the cover, the binding strip, and the book block spine, the cover can be provided on the transfer path into the pressing station with two parallel strip-shaped glue applications having a larger spacing than the thickness of the book block.\nAdvantageously, the compressed air supply device directed onto the binding strip parts is arranged oppositely to at least one conveying roll supporting the cover so that a counterpressure is built up which forces the respective parts against one another.\nExpediently, the compressed air supply device is connected with a compressed air source and can be controlled or adjusted with respect to the supplied amount of air and/or the speed of the compressed air.\nAdvantageously, the compressed air supply device is arranged so as to be adjustable in accordance to the thickness of the book block to be processed and is thus adaptable correspondingly.\nThe method according to the invention can be used preferably in connection with a back gluing machine or the like for the manufacture of books."}
{"text": "The following references are referenced in the specification and may be relevant to understanding the invention:    1. Gamble, J. and Milner, J. (1988). Evidence that immunological variants of p53 represent alternative protein conformations. Virology, 162, 452-458.    2. Harlow, E., Crawford, L. V., Pim, D. C. and Williamson, N. H. (1981). Monoclonal antibodies specific for simian virus 40 tumor antigens. J. Virol, 39, 861-869.    3. Yewdell, J. W., Gannon, J. V. and Lane, D. P. (1986). Monoclonal antibody analysis of p53 expression in normal and transformed cells. J. Virol. 59, 444-452.    4. Milner, J., Cook, A. and Sheldon, M. (1987). A new anti-p53 monoclonal antibody, previously reported to be directed against the large T antigen of simian virus-40. Oncogene, 1, 453-455.    5. Gannon, J. V., Greaves, R., Iggo, R. and Lane, D. P. (1990). Activating mutations in p53 produce common conformation effects. A monoclonal antibody specific for the mutant form. EMBO J. 9, 1595-1602.    6. Marasco, W. A. (1995) Intracellular antibodies (intrabodies) as research reagents and therapeutic molecules for gene therapy. Immunotechnology, 1, 1-19.    7. Stephen. C. W. and Lane, D. P. (1992). Mutant conformation of p53 precise epitope mapping using a filamentous phage epitope library. J Mol Biol. 225, 577-583.    8. Jannot. C. B. and Hynes, N. E. (1997). Biochem. Biophys. Res. Commun., 230, 242-246.    9. Cohen, P. A., Mani, J. C. and Lane, D. P. (1998). Characterization of a new intrabody directed against the N-terminal region of human p53. Oncogene, 17, 2445-2456.\nThe p53 gene encodes a protein which is 53 kD in size (hence the name). The p53 protein is situated in the nexus of intermingled pathways controlling cell proliferation, cell survival and differentiation. It can sense and integrate various external and internal stimuli, such as DNA damage, hypoxia, oxidative stress, deregulated oncogene expression and ribonucleotide depletion. In response to these stimuli, it may trigger cell cycle arrest, apoptosis, senescence, differentiation or antiangiogenesis. It appears that the p53 protein is involved in regulation of critical growth controlling checkpoints. It is capable of exercising a tumor suppressor function by preventing cells, that are in unfavorable environmental conditions or are carrying damaged DNA, from entering a cell cycle.\nMutations within the p53 gene have been found in more than 50% of all human cancers, rendering it the most frequently mutated single gene in human cancer known so far. Most of the mutations in the p53 gene discovered in cancers are of a missense type. They cause not just abrogation of the tumor suppressor function of wild-type p53, but often actively contribute to the tumor transformation function of mutant p53. Inactivation of p53 is one of the most common molecular events in cancer development.\nAs mentioned above, the p53 protein, which is both a regulator for cell proliferation and a suppressor of tumor development, can prevent the development of cancer by blocking the division of cells which have sustained DNA damage, or by triggering apoptosis. It has been proposed that the effect of mutant p53 on tumor progression is due to a dominant negative interaction of the wild type and mutant proteins. Thus, wild-type and mutant p53 proteins are respectively capable of contrasting suppressor and promoter effects on tumor development.\nThe ability of p53 to act both as a suppressor and a promoter of tumor development may reflect the ability of the protein to adopt different conformations. Mitogenic stimulation of primary T cells induces a change in the immunoreactivity of p53 and this may be due to a change in the tertiary structure of the p53 protein. To ascertain whether the ‘mutant’ conformation of wild-type p53 has physiological relevance, attempts have been made to associate it with a biochemical activity. Conformation-specific monoclonal antibodies, previously shown to discriminate between wild-type and mutant p53 proteins, have been used to demonstrate structural changes in wild-type p53 following sequence-specific binding to DNA (1). These studies suggested that wild-type p53 can physiologically adopt distinct conformations, which determine its DNA binding activity. Mutations that render p53 oncogenic may lock p53 into one of the few conformational states it physiologically adopts, rather than distort its tertiary structure.\nTwo important cellular factors appear necessary for this conversion: (1) binding of an ‘activating’ polypeptide that causes neutralization of the C-terminal negative regulatory domain, and (2) a highly reduced environment which can maintain p53 in an activated state. Given the clear association between the DNA binding activity and tumor suppressor functions of p53, these results imply that in many tumor cells there are high levels of mutant p53 that can potentially be activated to restore significant wild type function.\nA variety of monoclonal antibodies have been prepared against various epitopes of p53.\nA panel of anti-p53 mouse monoclonal antibodies has been used to characterize the immunoreactivities of the native p53-Ala35 and the mutant p53-Val 35 translated under various conditions. Two monoclonal antibodies, PAb421 and PAb248, were able to recognize discrete denaturation-stable epitopes on the p53-polypeptide (2,3). These antibodies immunoprecipitated both p53-Ala35 and p53-Val 35 translated at 30° C. and 37° C., and were also able to immunoprecipitate p53 from a variety of murine cell lines (4).\nThree additional monoclonal antibodies, PAb246, PAb 1620 and PAb240, were found to detect conformational changes in the p53 protein (5). Molecules in the mutant conformation are distinguished from wild-type molecules, inter alia, by the appearance of a new, normally cryptic epitope recognized by Pab240. This epitope was localized to residues 213 to 217 of the p53 protein and has the sequence RHSVV, preceded by F in human and mouse p53 (7). More than 90% of mutations found in p53 produce a conformational change in the p53 protein which results in the exposure of this epitope, which is otherwise hidden in the hydrophobic core of the molecule. This epitope will be referred to hereinafter as the common mutant epitope of p53.\nRecent advances in antibody engineering have allowed the genes encoding antibodies to be manipulated, so that antigen binding molecules can be expressed within mammalian cells in a controlled way (6). Application of gene technologies to antibody engineering has enabled the synthesis of single-chain fragment variable (scFv) antibodies that combine within a single polypeptide chain the light and heavy chain variable domains of an antibody molecule covalently joined by a predesigned peptide linker. The resultant scFv gene can be expressed in bacterial expression systems such as E. coli. Bundled in the “gene display package” single-chain antibodies displayed at the surface of filamentous phages of the M13 family provided the possibility to create antibody libraries both from various living sources and products of diversification of a single scFv molecule. Antibodies with the desired specificity can be isolated from such libraries employing effective selection techniques (biopanning) in which the antigen is immobilized on a solid support.\nThe ability to create scFv antibodies, when combined with their stable expression in precise intracellular locations in mammalian cells, has enabled the creation of a powerful new family of antibody molecules for basic research or gene therapy. These intracellular antibodies (intrabodies) can be used to modulate cellular physiology and metabolism through a variety of mechanisms, including the blocking, stabilizing or mimicking of protein-protein interactions, by altering enzyme function, or by diverting proteins from their usual intracellular compartments. Intrabodies can be directed to the relevant cellular compartments by modifying the genes that encode them to specify N- or C-terminal polypeptide extensions for providing intracellular-trafficking signals. This approach has been described for a number of different antigens, including several HIV proteins.\nSeveral scFv antibodies against p53 have been previously described. The scFv-421 antibody recognizes a C-terminal epitope of the protein (8). When expressed in vitro in the cytoplasm or in the nucleus of COS-1 cells, it was found to be non-functional and prone to rapid degradation. An important determinant of correct antibody folding is the formation of intra-chain disulfide bonds in the variable regions; possibly the reducing environment of the cytosol may lead to a decrease in the stability of the scFv (9). Nevertheless, some scFv antibodies have been expressed in the cytoplasm and shown to have biological effects, indicating that other features, such as the primary sequence of the antibody and/or its specific cellular location may be important for their proper function.\nThe scFv DO-i antibody was found to recognize an N-terminal epitope of human p53 (9). The DO-1 scFv was targeted to the cytoplasm and to the nucleus of mammalian cells. Interestingly, insertion of the CK domain into scFv to create a scFvCK fusion protein, led to a dramatic increase in the level of intracellular expression. However, in other studies in which CK fusions to scFv were made, the effects on stability and expression were much less marked. It is clear that each scFv is a particular and individual case (9)."}
{"text": "A polymer material having a carbon-carbon double bond in the molecule thereof such as EPDM, NR and NBR is susceptible to degradation by sunlight (particularly ultraviolet rays) or ozone, so that a problem of weatherability may occur depending on the application. In addition, such a polymer material may be dissolved in a specific organic solvent or the like, and cannot be applied also to a site that may be contacted with such an organic solvent or the like. Therefore, as one countermeasure for extending the application of such a polymer material, there is considered a modification (chemical modification) by a nitrile oxide.\nOn the other hand, it has been known from long ago that a nitrile oxide which is one of representative 1,3-dipoles, is reacted with a monomer of alkene or alkyne, a monomer of nitrile compounds and the like under a moderate condition. However, a nitrile oxide is, because of its high reactivity, gradually dimerized in a non-polar solvent to be converted to furoxane, so that the isolation and the preservation of a nitrile oxide is difficult. Therefore, in order to perform a modification by a nitrile oxide, a method of generating a nitrile oxide in the reaction system by reacting a base such as triethylamine with halogenated aldoxime that is a precursor of a nitrile oxide becomes general.\nHowever, by this method, when it was attempted to react a polymer material such as NBR and PAN with a nitrile oxide, the reaction hardly progressed, so that a modified polymer material modified by a nitrile oxide cannot be obtained.\nHere, Japanese Patent Application Publication JP-A-11-180943 discloses a synthetic method of benzonitrile oxide derivatives such as mesitylenedinitrile oxide and mesitylenemononitrile oxide, but discloses no polymer material modified by a benzonitrile oxide derivative.\nIn addition, Japanese Patent Application Publication JP-A-2008-163232 discloses that by using a compound having a heterocycle such as 4-(2-oxazolyl)-phenyl-N-phenylnitrone and 4-(2-oxazolyl)-phenyl-N-methylnitrone, a diene-based copolymer rubber is kneaded and reacted in cyclohexane, but discloses no polymer material modified by an aromatic nitrile oxide having a substituent at an ortho-position."}
{"text": "Fluid filters are widely known and used. After a period of time and usage, many fluid filters undergo wear, and eventually must have their filter element(s) serviced and/or replaced. Many fluid filters also have multiple parts that must be removed in order to access the filter element(s).\nHowever, many fluid filters are structured to be relatively disposed in low tolerance areas having rather tight fits, so that such fluid filters occupy a minimal amount of space. As a result, servicing such fluid filters and replacing their filter element(s) can be difficult. For example, servicing the fluid filter has involved moving or removing surrounding components of an engine and other parts of the fluid filter to gain access to the filter element(s). Furthermore, servicing has been known to be messy due to such spatial constraints and dripping has been known to occur when removing the fluid filter to replace the filter element(s). Thus, improvements can still be made to existing fluid filters and to existing fluid filter servicing processes."}
{"text": "Field of the Invention\nThe present invention relates, in general, to a haptic apparatus and, more particularly, to a locally vibrating haptic apparatus and a method for locally vibrating the haptic apparatus, capable of vibrating a desired position of the haptic apparatus or adjusting a vibrating position by improving arrangement of vibrators and a frequency control method, and to a haptic display apparatus and a vibration panel having a plurality of excitation points, capable of isolating vibration from a housing to vibrate a desired position of a display part or adjust a vibrating position, vibrating a desired position of the display part or adjusting a vibrating position.\nDescription of Related Art\nA display apparatus having a touch window is being widely spread as an improvement in sensitive manipulation of an interface is recently demanded. The touch window is a pointer input apparatus that is similar to a mouse, and is operated in conjunction with an image display apparatus such as a light emitting diode (LED) or a liquid crystal display (LCD).\nThe touch window is an apparatus that directly points a position by hand, or moves a pointer and inputs a desired instruction through a specific movement while being in contact with the touch window using an exclusive input tool such as stylus.\nSuch a touch window is configured to be intuitively used by stimulating a user sense, mainly, sight and hearing. Recently, an apparatus employing a haptic technology using a sense of touch is a growing trend.\nThe haptic technology is being widely applied to an apparatus having a display, such as a mobile apparatus, a monitor or a television.\nGenerally, a touch screen that is an LCD for forming an image and inputting a command through touch is provided on a front of the mobile apparatus. In order to execute a command displayed on the touch screen, a pointer is moved around the screen by a user manipulation. When the pointer reaches a desired menu or position, the command is executed by tapping at the desired menu or position with a finger.\nIn order to improve the sensitive manipulation of the touch screen, a vibrator is provided to generate vibration. The vibrator is generally received in a main body of the mobile apparatus to transmit vibration to the main body or to the window.\nIf a user inputs a command by touching a specific portion using a finger or a stylus, vibration is generated by the vibrator through which rated voltage flows, thus feeling vibration.\nThe vibrator may use a vibration motor or a linear actuator. Such a vibrator transmits vibration to the screen or the main body.\nHowever, such a vibrator serves to simply transmit vibration. Thus, technical development for delicate and precise vibration control is required to achieve an original function of the haptic apparatus.\nFurther, a different magnitude of vibration may be frequently generated for each portion of the screen according to a location of the vibrator. The vibrator arranged as such may generate a larger magnitude of vibration in the main body than in the screen."}
{"text": "The present inventor has appreciated that while children and adults are fascinated by the flow of water that games and, particularly, competitive games typically do not utilize a player's ability to control the flow of water in a game device.\nPreviously known game devices frequently do not place the game players under time pressures to perform and, particularly, under time pressures which are variable depending upon the play and are not merely directed to a counting of time."}
{"text": "1. Field of the Invention\nThe present invention relates generally to methods of manufacturing semiconductor devices, and more specifically, to a method of forming a well region having a prescribed impurity concentration distribution in a region positioned under an isolation insulating film.\n2. Description Of The Background Art\nA typical example of a semiconductor device having MOS transistors formed on a main surface of wells formed on a semiconductor substrate is a complementary MOS device (hereinafter referred to as CMOS device). A CMOS device is characterized in that an n channel MOS transistor and a p channel MOS transistor are formed in a single semiconductor substrate. The power consumption of a CMOS device is advantageously very small, because DC current flowing between power supply terminals is very small. Now, the structure of a conventional CMOS device will be described as the background of the invention in conjunction with the drawings.\nFIG. 43 is partly sectional view showing the structure of a conventional CMOS device. An n well 5 and a p well 6 having different conductivity types from each other are formed in a surface region of a p type silicon substrate 1. An isolation oxide film (field oxide film) 2 for element isolation is formed in a prescribed region on surfaces of n well 5 and p well 6. A p channel MOS transistor 50 is formed on a surface of n well 5. An n channel MOS transistor 60 is formed on a surface of p well 6. P channel MOS transistor 50 has a gate electrode 8, and a pair of p.sup.+ impurity regions 9a, 9b spaced apart from each other. N channel MOS transistor 60 has a gate electrode 8, and a pair of n.sup.+ impurity regions 10a, 10b spaced apart from each other. Note that although only one transistor is depicted in each of n well 5 and p well 6 in FIG. 43, a plurality of MOS transistors and other functional elements are formed in practice.\nThe CMOS structure is however susceptible to a latch up phenomenon in which excessive current flows across power supply terminals to break down elements. The latch up herein indicates a phenomenon in which pnp and npn parasitic bipolar transistors are formed in a CMOS device to constitute a pnp thyristor between power supply potential (V.sub.DD) and ground potential GND (V.sub.SS), and therefore current continues to flow between V.sub.DD and GND once extraneous noise is applied, resulting in breakdown.\nFIG. 44 is a cross sectional view schematically showing a parasitic thyristor formed in a CMOS device the same as one described in conjunction with FIG. 43. In FIG. 44, if the impurity concentrations of n well 5 and p well 6 are low, voltage drop down (voltage drop down corresponding to resistors Rn, Rp) increases when some surge is applied and current flows across these well regions. Thus, the emitter-base regions of parasitic pnp bipolar transistor Q1 and parasitic npn bipolar transistor Q2 are biased. As a result, these parasitic transistors operate and the latch up phenomenon described above is likely to occur.\nA so-called retrograde well structure with an increased impurity concentration at the bottom of well is employed for the purpose of improving the resistance against the latch up phenomenon. The retrograde well is usually formed by implanting impurity ions with high energy into a semiconductor substrate.\nA method of forming a retrograde well by means of implanting impurity ions with a high energy into a semiconductor substrate is, for example, disclosed in John Yuan-Tai Chen, \"Quadruple-Well CMOS for VLSI Technology, \" IEEE Transactions on Electron Devices, vol. ED-31, No. 7, July 1984 and U.S. Pat. No. 4,633,289. FIGS. 45-49 are cross sectional views showing steps in the manufacture of a conventional retrograde well structure.\nReferring to FIG. 45, a thick isolation oxide film 22 is selectively formed on a surface of a p type silicon substrate 1 so as to separate element formation regions. Isolation oxide film 22 is formed by means of LOCOS (Local Oxidation of Silicon). The LOCOS is a process of thermally oxidizing an underlying oxide film 3 having only a prescribed surface region exposed by a patterned nitride film. The formation of isolation oxide film 22 defines an active region in which MOS transistor and the like are to be formed. Underlying oxide film 3 is formed on the element formation region.\nThen, as illustrated in FIG. 46, a region to form the p well is covered with patterned resist 41. Using patterned resist 41 as mask, phosphorus ions (P.sup.+) are implanted with high energy into silicon substrate 1 a number of times, changing the energy and the implantation amount. Thus, n type retrograde wells 51, are formed to have a first n type concentration peak position 51a at a deep position in silicon substrate 1 and a second n type impurity concentration peak position 51b under isolation oxide film 22. At the time, low energy ion implantation is not performed to prevent increase of the impurity concentration of the n type well in a shallow region up to about 1000 .ANG. in depth from the surface of silicon substrate 1.\nAs illustrated in FIG. 47, after removal of patterned resist 41, the n well region is covered with patterned resist 42. Using patterned resist 42 as mask, boron ions (B.sup.+) as p type impurity ions are implanted with high energy a number of times into silicon substrate 1, changing the energy and the implantation amount. Thus, p type retrograde wells 61, are formed to have a first p type impurity concentration peak position 61a at a deep position in silicon substrate 1, and a second p type impurity concentration peak position 61b under isolation oxide film 22. As in the above case, low energy ion implantation is not performed to prevent increase of the impurity concentration of the p type well in the vicinity of the surface of silicon substrate 1. Thereafter, as illustrated in FIG. 48, patterned resist 42 is removed away.\nFinally as illustrated in FIG. 49, after removal of underlying oxide film 3, a gate oxide film 7 is formed again in the same region. Gate electrode 8 is formed on gate oxide film 7. A pair of p.sup.+ impurity regions 9a, 9b are formed spaced apart from each other with gate electrode 8 therebetween in n type retrograde well region 51. A pair of n.sup.+ impurity regions 10a, 10b are formed spaced apart from each other with gate electrode 8 therebetween in p type retrograde well region 61. Thus, a p channel MOS transistor 50 is formed in n type retrograde well region 51, while n channel MOS transistor 60 is formed in p type retrograde well region 61.\nIn the retrograde well regions thus formed, the region having first impurity concentration peak positions 51a, 61a formed at the deep positions of silicon substrate 1 is effective in preventing the latch up phenomenon. The region having second impurity concentration peak positions 51b, 61b functions as a channel stop region for element isolation.\nAs the size of elements such as a MOS transistor formed in a semiconductor substrate is reduced, the thickness of an isolation oxide film is reduced. More specifically, as the element isolation width is reduced from the order of micron to the order of submicron, the thickness of the isolation oxide film is reduced from about 5000 .ANG. to about 3000 .ANG..\nThe reason for the thickness of the isolation oxide film being reduced will be described. As the size of elements such as MOS transistor is reduced, the element isolation width should be reduced. More specifically, the element isolation width should be reduced to the order of submicron less than 1 .mu.m and further less than 0.5 .mu.m. In order to achieve such an element isolation width, a field oxidation treatment must be performed after patterning a nitride film to have an opening width less than 1 m.\nNow, the relation between the thicknesses of isolation oxide films and the opening widths of nitride films when each isolation oxide film is formed under the same conditions will be described. The relation between the thickness of isolation oxide film and the opening width of nitride film under the same treatment conditions is disclosed in \"Oxidation Rate Reduction in the Submicrometer LOCOS Process\", IEEE TRANSACTIONS, Vol. ED-34, No. 11, Nov. pp. 2255-2259 1987.\nAccording to this document, when a field oxidation treatment was performed using a nitride film having an opening width less than 1 .mu.m, it was demonstrated the thickness of the isolation oxide film was reduced as the opening width was reduced. This means that the thickness of the isolation oxide film is reduced as the element isolation width is reduced under the same oxidation treatment conditions. However, even in the case of a nitride film having a small opening width, the thickness of an isolation oxide film can be increased by prolonging time for treatment. Thus prolonging the time for treatment is however disadvantageous in view of reducing size, because bird's beaks are increased accordingly. In view of the foregoing, it is pointed that thickness of the isolation oxide film is reduced as the element isolation width is reduced.\nFIG. 50 is a cross sectional view showing the state of the retrograde well structure shown in FIG. 49 in which the isolation oxide film is reduced. As illustrated in FIG. 50, isolation oxide film 23 is reduced in thickness as compared to isolation oxide film 22 in FIG. 49. Accordingly, second impurity concentration peak positions 51b, 61b forming the retrograde wells are formed at the positions closer to the surface of silicon substrate 1. This is because impurity ions are implanted into silicon substrate 1 through thin isolation oxide film 23, and second impurity concentration peak positions 51b, 61b formed as a channel stop region for element isolation are formed at the positions closer to the surface of silicon substrate 1. As a result, the impurity is diffused by thermal treatment, etc. in succeeding steps and is likely to reach the vicinity of the surface of silicon substrate 1. This adversely affects the transistor characteristics, especially substrate biasing effect on the surface of silicon substrate 1.\nThe substrate biasing effect herein indicates the effect in which the threshold voltage of an MOS transistor is changed by applying voltage to the silicon substrate in which the MOS transistor is formed. More specifically, the threshold voltage is proportional to the square root of the substrate bias voltage. The constant of proportion is defined as a substrate effect constant. The adverse effect upon the above-described substrate biasing effect indicates increase of the substrate effect constant.\nFIG. 51 is a graph showing the relation between threshold voltage V.sub.th and substrate bias voltage V.sub.BS. When the thickness of an isolation oxide film is 5000 .ANG. and an impurity concentration peak position is present under the isolation oxide film, as illustrated in FIG. 51, the substrate effect constant K will be 0.2. In this case, even if substrate bias voltage V.sub.BS changes from 0 V to -5 V, the degree of change of threshold voltage Vth is relatively small. If, however, thickness of the isolation oxide film is reduced to about 2000 .ANG. as reduction of the size of elements proceeds and a retrograde well is formed to have an impurity concentration peak position under the isolation oxide film, an impurity is implanted to a shallow position from the substrate surface. Accordingly, the impurity diffuses due to thermal treatment in subsequent steps, resulting in increase of the impurity concentration in the vicinity of the substrate surface. As a result, the substrate effect constant K becomes as large as 0.5. Thus, slight change of a substrate bias voltage V.sub.BS results in great change of threshold voltage V.sub.th.\nIn this case, a bias voltage imposed state is encountered due to extraneous noise, even if bias voltage is not applied to the substrate. Therefore in a CMOS device having a retrograde well structure, when the isolation oxide film is reduced in thickness with reduction of the size of transistors, the threshold voltage of an M0S transistor is likely to change if bias voltage is supplied to the substrate or a bias voltage imposed state is attained due to some extraneous causes. Thus, as the thickness of the isolation oxide film is reduced, the impurity performing the retrograde well structure adversely affects the characteristic of the MOS transistors formed on the substrate surface.\nA disadvantage encountered when a retrograde well structure is formed in a shallow portion of the above-described element formation region will be described in more detail in conjunction with FIGS. 52-55. These FIGS. 52-54 are cross sectional views showing manufacturing steps especially in view of the retrograde well structure. FIG. 55 is a view for use in illustration of how the substrate biasing effect is influenced when a retrograde well structure is formed in a relatively shallow portion in an element formation region.\nReferring to FIG. 52, an isolation oxide film 102 is selectively formed in a main surface of p type semiconductor substrate 101. A first impurity concentration peak position 105a is formed by implanting impurity ion with a high energy through isolation oxide film 102. Then, referring to FIG. 53, an impurity is implanted into a second impurity concentration peak position 105b by implanting the impurity with energy the same as impurity implantation in the vicinity of the bottom surface of isolation oxide film 102. After the impurity is thus implanted, a prescribed thermal treatment is performed to form a second impurity concentration peak position 105b extending from the vicinity of the bottom surface of isolation oxide film 102 to the bottom of element formation region.\nSince second impurity concentration peak position 105b is thus formed, the depth D of second impurity concentration peak position 105b under the element formation region is determined by the thickness t3 of isolation oxide film 102. More specifically, if the thickness t3 of isolation oxide film 102 takes a small value, second impurity concentration peak position 105b will be formed at a relatively shallow position in the element formation region. As a result, the substrate biasing effect in the element formation region is increased.\nThe reason for the increase of the substrate biasing effect by the formation of the second impurity concentration peak position at such a shallow position from the semiconductor substrate surface will be described. The threshold voltage Vth of a transistor is given by the following equation in Mitsumasa Koyanagi, Electronic Material Series Submicron Device I, MARUZEN KABUSHIKI KAISHA, pp. 4-8. ##EQU1## where V.sub.FB represents flat band voltage, .PHI..sub.F Fermi potential of substrate, V.sub.BS substrate bias voltage, and C.sub.OX oxide film capacitance. In the above equation, (2.multidot..epsilon..sub.S .multidot..epsilon..sub.0 .multidot.q.multidot.N.sub.A).sup./C.sub.OX is a coefficient called substrate effect constant. The substrate effect constant (K) is a coefficient representing the degree of modulation of V.sub.th by the substrate biasing. The dependence of the substrate biasing effect upon the impurity concentration of the substrate is due to acceptor concentration N.sub.A included in the above-described substrate effect constant. More specifically, increase of acceptor concentration N.sub.A in the region in which an inversion layer and a depletion layer are formed increases the value of threshold voltage V.sub.th.\nFIG. 55 shows an MOS transistor having a conventional retrograde well structure. Referring to FIG. 55, a source region 114 and a drain region 113 are formed to define a channel region in an element formation region. A gate electrode 111 is formed on the channel region with a gate insulating film 112 therebetween. Voltage V.sub.D is applied to the drain region of the MOS transistor, and voltage V.sub.G is applied to gate electrode 111. Source region 113 is grounded.\nSupplying prescribed voltages to drain region 113 and gate electrode 111 forms a depletion layer 110. In the vicinity of the region in which depletion layer 110 is formed, however, as illustrated in FIG. 55, the second impurity concentration peak position 105b of retrograde well structure 105 will be formed. This restricts depletion layer 110 from expanding, and as illustrated in FIG. 55, for example, the depletion region will be formed into a narrow depletion layer 110a. Thus, the depletion layer capacitance is increased, whereby the substrate biasing effect is increased accordingly.\nMore specifically, the formation of the impurity concentration peak position at a shallow position in the element formation region increases the accepter concentration in the vicinity of the region to form the depletion layer, thus increasing the substrate effect constant. The substrate effect constant takes the value of K shown in FIG. 51, and as the value of the substrate effect constant is larger, change of the threshold voltage V.sub.th is increased by supplying the substrate bias voltage. More specifically, the substrate biasing effect is increased.\nMeanwhile, one method of forming a channel cut layer under an isolation oxide film is to previously implant an impurity for forming the channel cut layer into the isolation oxide film formation region and then to form the isolation oxide film. According to this method, the problem associated with the substrate biasing effect described above can be alleviated, but the problem of narrow channel effect is encountered. This will be described in conjunction with FIGS. 56-62. FIGS. 56-60 are cross sectional views showing manufacturing steps for forming an element isolation structure by which a channel cut layer is previously implanted with an impurity, and then an isolation oxide film is formed.\nReferring to FIG. 56, an impurity is implanted to form a well in a p type semiconductor substrate 101, and a first impurity concentration peak 105a iis formed at a deep position in p type silicon substrate 101. Then, an oxide film 102a and a nitride film 103 are sequentially formed on a surface of p type silicon substrate 101. As illustrated in FIG. 57, resist 106 which is patterned into a prescribed form is formed on nitride film 103, and nitride film 103 is patterned into a prescribed form using resist 106 as mask.\nReferring to FIG. 58, an impurity 104a is implanted to form a channel cut layer using the above-described resist 106 as mask. Then, referring to FIG. 59, after removal of resist 106, an isolation oxide film 102 is selectively formed by means of thermal oxidation treatment. Thus, channel cut layer 104 is formed under isolation oxide film 102. Subsequently, as illustrated in FIG. 60, nitride film 103 is removed away.\nNow, referring to FIGS. 61 and 62, the problem of the narrow channel effect associated with the above-described structure will be described. FIG. 61 is a cross sectional view showing an MOS transistor having the above-described element isolation structure. FIG. 62 is a plan view showing the MOS transistor including the above-described element isolation structure. A cross section taken along line A--A in FIG. 62 corresponds to FIG. 61.\nReferring to FIG. 61, an element formation region 109 is defined by isolation oxide film 102. channel cut layer 104 is formed to extend to the vicinity of the surface of element formation region 109. In this case, a substantial element formation region width W2 is smaller than a desired element formation region width W3 by the function of channel cut layer 104. More specifically, channel cut layer 104 is formed to enter element formation region 109 by the width of 2S, and the width of element formation region is reduced by the amount. The narrow channel effect is thus enhanced. When viewed two-dimensionally, this is as illustrated in FIG. 62. More specifically, channel cut layer 104 is formed around element formation region 109 to stick out from isolation oxide film 102. This narrows the substantial width of element formation region 109. The narrow channel effect thus become remarkable. In contrast, in the above-described retrograde well structure, the narrow channel effect described above is hardly encountered."}
{"text": "In the present day, the SHVC encoding technology (scalable expansion of HEVC) is known that enables expansion to super-resolution videos such as 4K resolution videos at a low bitrate using HDTV video transmission. Herein, HDTV stands for High-definition television, SHVC stands for Scalable High-efficiency Video Coding, and HEVC stands for High-efficiency Video Coding.\nIn the SHVC encoding technology, in addition to performing inter-screen prediction and in-screen prediction used in HEVC, inter-layer prediction meant for predicting 4K images from among HDTV images is also performed; and the subjective image quality is improved by holding down the volume of data assigned to such image areas which have movement-induced distortion and by assigning a greater volume of data to such in-focus image areas which have less movement. Moreover, temporal-spatial image processing is performed so that the difficult-to-compress noise components are removed in advance from the 4K images. That enables achieving control on the decline in the feel of resolution at a low bit rate. For that reason, even in the case in which the additional bit rate is set to 10 Mbps, it becomes possible to maintain the feel of resolution of the 4K images. In order to distribute videos of 4K resolution, it is necessary to have the bit rate of, for example, about 30 Mbps. However, as a result of implementing the HEVC encoding technology, the additional transmission bandwidth that is required can be reduced to about 10 Mbps that is one-third of 30 Mbps.\nIn a narrow-bandwidth video transmission system in which SHVC is used, reference images are shared in advance between the transmission side and the reception side. The transmission side transmits encoding data representing only the difference with the reference images. The reception side decodes the received encoding data using the reference images shared in advance. Since the encoding data represents only the difference with the reference images, it becomes possible to reduce the volume of data and to perform video streaming having excellent bandwidth usage efficiency."}
{"text": "1. Field of the Invention\nThe present invention relates to a power conversion controller, especially to a fixed-on-time controller for discontinuous mode PFC (Power Factor Correction) power conversion.\n2. Description of the Related Art\nFixed-on-time power conversion is capable of achieving good PFC for discontinuous mode power conversion when the line input is a full-wave rectified voltage. The reason is as follows:\nLet the line input be VIN(t)=|VA sin(ωt)|, the fixed on time=tON, the inductance=L, and the input current at the end of an on period=IINPEAK(tX), then IINPEAK(tx)=(VIN(tX)/L)tON=(|VA sin(ωtX)|/L)tON, wherein tON is much smaller than 2π/ω. As can be seen in the above equation, the input current will follow the line input at every on period and a good power factor correction is therefore achieved.\nTo get insight into the principle of the fixed-on-time discontinuous mode PFC power conversion, please refer to FIG. 1, which illustrates a block diagram of a discontinuous mode PFC power converter using a prior art fixed-on-time controller. As illustrated in FIG. 1, the discontinuous mode PFC power converter includes an input rectifier circuit 101, a transformer circuit 102, an output rectifier circuit 103, a load 104, a resistor 105, and a fixed-on-time controller 110.\nThe input rectifier circuit 101 is used to generate a full wave rectified voltage VFULL—WAVE according to an AC power VAC.\nThe transformer circuit 102 is used to transfer the power from the full wave rectified voltage VFULL—WAVE to the output rectifier circuit 103, under the control of a driving signal VDRV of the fixed-on-time controller 110.\nThe output rectifier circuit 103 is used to provide a DC voltage VO to the load 104, and the resistor 105 is used to provide a feedback signal VFB for the fixed-on-time controller 110.\nThe fixed-on-time controller 110 includes an error amplifier 111, a constant current source 112, a capacitor 113, a switch 114, a comparator 115, and a fixed-on-time driver circuit 116.\nThe error amplifier 111 has a negative input end coupled to the feedback signal VFB, a positive input end coupled to a reference voltage VREF, and an output end for providing a threshold signal VCOMP.\nThe constant current source 112, the capacitor 113, and the switch 114 are used to generate a saw signal VSW.\nThe comparator 115 has a negative input end coupled to the threshold signal VCOMP, a positive input end coupled to the saw signal VSW, and an output end for providing a turn-off signal VOFF.\nThe fixed-on-time driver circuit 116 is used to provide the driving signal VDRV for the transformer circuit 102 and a reset signal RESET for the switch 114, according to the turn-off signal VOFF from the comparator 115 and a sensing signal VAUX from the transformer circuit 102, wherein the sensing signal VAUX is used to indicate the end of an inductor current discharging period of the transformer circuit 102, and the active time point of the reset signal RESET follows that of the turn-off signal VOFF.\nWhen in operation, the voltage of the driving signal VDRV will arise from a low level to a high level after the sensing signal VAUX becomes active, and fall from a high level to a low level after the turn-off signal VOFF becomes active. Besides, the period the driving signal VDRV remains at a high level—the on time of the transformer circuit 102—will be fixed to a value by the feedback control mechanism of the power conversion to transfer a specific amount of energy per cycle from the AC power VAC to the load 104, to regulate VFB at VREF.\nHowever, as the period the driving signal VDRV remains at a high level—a fixed on time of the transformer circuit 102 corresponding to a load value of the load 104 and a line voltage of the AC power VAC—is ended when the saw signal VSW reaches the threshold signal VCOMP, the threshold signal VCOMP will exhibit a large level shift to change the period of the fixed on time from a short/long value to a long/short value when the load value of the load 104 or the line voltage of the AC power VAC changes drastically.\nPlease refer to FIG. 2, which illustrates the waveforms of major signals in the prior art fixed-on-time controller 110 of FIG. 1 corresponding to a low line and a high line of the AC power VAC respectively. As illustrated in FIG. 2, VDRV, LOW—LINE (the driving signal VDRV generated corresponding to a low line of the AC power VAC) has a fixed on time tON1, and VDRV, HIGH—LINE (the driving signal VDRV generated corresponding to a high line of the AC power VAC) has a fixed on time tON2, wherein tON1 is longer than tON2 such that same power is delivered to the load 104. As the fixed on time is ended when the saw signal VSW—of which the ramping up slope is fixed by the constant current source 112—reaches the threshold signal VCOMP, the fixed on time is then controlled by the level of the threshold signal VCOMP. As such, a higher VCOMP, LOW—LINE (the threshold signal VCOMP corresponding to the low line of the AC power VAC) is generated to allow a higher VSW, LOW—LINE (the saw signal VSAW corresponding to the low line of the AC power VAC) and therefore the longer tON1; and a lower VCOMP, HIGH—LINE (the threshold signal VCOMP corresponding to the high line of the AC power VAC) is generated to allow a lower VSW, HIGH—LINE (the saw signal VSAW corresponding to the high line of the AC power VAC) and therefore the shorter tON2, wherein the VCOMP, LOW—LINE is higher than the VCOMP, HIGH—LINE by ΔV1.\nThat is, ΔV1 of the prior art fixed-on-time controller 110 can be a large value when the load value of the load 104 or the line voltage of the AC power VAC changes drastically. However, large ΔV1 is adverse to the design of the error amplifier 111. To minimize ΔV1, one solution is to change the capacitance of the capacitor 113 to cope with the amplitude variation of the AC power VAC or the load value variation of the load 104. However, it is bothersome to change the capacitor 113 whenever the AC power VAC or the load 104 is changed, and besides, the solution is adverse to the integration of the capacitor 113 into the fixed-on-time controller 110.\nIn view of these problems, the present invention proposes a novel fixed-on-time controller for discontinuous mode PFC power conversion."}
{"text": "Various communication standards, such as digital subscriber line (xDSL), very-high-bit-rate digital subscriber line 2 (VDSL2), G.hn, and G.fast, have been proposed or developed to provide high-speed data transmission from the service provider (e.g., a central office) to a customer premise over the existing twisted-pair copper wiring conventionally used for telephone service. Such technologies leverage modem technology to increase the data transfer bandwidth of the twisted-pair copper wiring. Typically, modems are provided on the ends of the subscriber line copper wiring to communicate between the central office and the customer premise. The manner in which the two modems communicate is established by the particular standard governing the communication. Because the existing telephone wire is used, the data signals are typically transferred out-of band with the voice band signals. Because different frequencies are used for the voice band and the data band, voice and data information can be concurrently transferred over the twisted-pair copper line.\nService providers have increased data bandwidth by installing fiber optic cabling between the central office (CO) and a distribution point unit (DPU) closer to the customers. A particular DPU may interface with a bundle of twisted pairs to service a relatively small number of customer premise connections. This approach shortens the length of the copper pair between the CO interface at the DPU and the customer, thereby allowing increased data rates. Thus the DPU will provided telephony and/or data to one or more customer premises equipment (CPE)\nOne difficulty arising from an optical connection between the central office and the DPU lies in the inability to provide a source of power for the DPU. Due to the remoteness of the DPU with respect to the central office, a local power supply is often unavailable or expensive to install.\nPowering for a DPU may be provided by reverse power feeding, wherein power is supplied to the DPU from the various CPEs for which telephony and/or data services are provided from the DPU. A standard for reverse power feeding is being standardized by ETSI and the Broadband World Forum. In such an embodiment, a power supply in the DPU may combine power contributions from multiple CPEs to power a main distribution unit (MDU) that handles the voice and data communication. This arrangement is referred to as a reverse power system, since the CPEs are the power sourcing equipment (PSE) and the DPU, particularly the MDU, is the powered device (PD).\nThe CPE PSE thus injects power across the copper pair. In order to use a plain old telephony service (POTS) type analog telephone, a POTS adapter is to be installed between the copper pair and the POTS telephone. FIG. 1 illustrates a high level block diagram of such an arrangement, comprising a DPU 20 and a CPE 30 connected by a copper pair 40, with certain standard ETSI symbols shown. DPU 20 comprise a data provision module 50, typically in optical communication with a CO; a DC extraction circuitry 60; and a distribution point power supply 65. Data provision module 50 and DC extraction circuitry 60 are each connected to copper pair 40 past the U-O reference point. The U-O reference point is defined as the reference point at the DPU containing both DC power and service data. Distribution point power supply 65 is arranged to convert power received from DC extraction circuitry 60 to an appropriate power for data provision module 50, and other devices located within DPU 20.\nCPE 30 comprises a PSE 70, a power splitter 80, a service splitter 90, a POTSA-D adaptor 100 and a POTS telephone 110. PSE 70 is connected to power splitter 80 across the U-R2P reference point, defined as the reference point at CPE 30 containing the injected DC power. Power splitter 80 is connected to service splitter 90, which provides service and optionally analog phone service for CPE 30. Power splitter 80 is additionally connected to copper pair 40 across the U-R reference point, defined as the reference point at CPE 30 containing both DC power and service data. POTSA-D 100 is connected to copper pair 40 between power splitter 80 and the U-R reference point. POTSA-D 100 is an adapter that can be attached to one or more POTS telephone 110 in CPE 30. POTSA-D 100 is arranged to perform the following functions: translate the signals from the upstream DC and low frequency POTS signaling from the POTS telephone 110 into an in-band or out-of-band signaling system; translate the signals from the downstream in-band or out-of-band signaling system into POTS signaling towards the POTS telephone 110; and provide sufficient current, with a current limit, and DC voltage to supply POTS telephone 110.\nIn the event that POTS telephone 110 is mistakenly connected to the reverse power feeding network without POTSA-D 100, i.e. POTS telephone 110 is mistakenly connected directly to the in-premises wiring at the potential of reference point U-R, when POTS telephone 110 goes off-hook it draws a large current from PSE 70. A standard POTS telephone 110 assumes that the CO supplies a current limited power, and thus is not designed with an inherent current limiter. As indicated above, one of the functions of POTSA-D 100 is thus to limit the current, however in the absence of POTSA-D 100, POTS telephone 110 may draw very large amounts of current from PSE 70, which may result in an unsafe condition, such as a fire.\nAn additional issue is one of identification by PSE 70 of the proper conditions over which to supply power. It is presently contemplated by the inventors that PSE 70 should provide for detection of DC extraction circuitry 60 prior to supplying power. As such, it is required that POTSA-D 100, which ultimately may also draw power from PSE 70, not interfere with such a detection process."}
{"text": "This invention relates to microelectromechanical systems (MEMS).\nMicroelectromechanical systems (MEMS) include micromechanical elements, actuators, and electronics. Micromechanical elements may be created using deposition, etching, and or other micromachining processes that etch away parts of substrates and/or deposited material layers or that add layers to form electrical and electromechanical devices. One type of MEMS device is called an interferometric modulator. As used herein, the term interferometric modulator or interferometric light modulator refers to a device that selectively absorbs and/or reflects light using the principles of optical interference. In certain embodiments, an interferometric modulator may comprise a pair of conductive plates, one or both of which may be transparent and/or reflective in whole or part and capable of relative motion upon application of an appropriate electrical signal. In a particular embodiment, one plate may comprise a stationary layer deposited on a substrate and the other plate may comprise a metallic membrane separated from the stationary layer by an air gap. As described herein in more detail, the position of one plate in relation to another can change the optical interference of light incident on the interferometric modulator. Such devices have a wide range of applications, and it would be beneficial in the art to utilize and/or modify the characteristics of these types of devices so that their features can be exploited in improving existing products and creating new products that have not yet been developed."}
{"text": "There is some demand for high-capacity and high-energy density secondary batteries as power supplies for mobile devices and electric vehicles. Secondary batteries degrade by, for example, being repeatedly charged and discharged. As secondary batteries degrade, the capacity of discharge from a fully charged state (hereinafter, simply referred to as the “capacity”) decreases. As a result, in the case of electric vehicles, the distance that can be traveled per charge becomes shorter. Accordingly, to make it possible to treat degradation of a secondary battery, for example, by replacing the secondary battery, at an appropriate time, there is some need to develop a method for sensing secondary battery degradation in a simplified and correct manner.\nIn view of this, there has been proposed a conventional determination of secondary battery degradation on the basis of an open-circuit voltage of a secondary battery and a charge capacity obtained by current integration (see Patent Document 1). Here, a capacity change is calculated from a pre-charge open-circuit voltage and a post-charge open-circuit voltage of the secondary battery, and the calculated capacity change is compared with a charge capacity obtained by current integration."}
{"text": "Distance Vector Multicast Routing Protocol (DVMRP) is a widely deployed multicast routing protocol. When used with other multicast routing protocols such as Protocol Independent Multicast-Sparse Mode (PIM-SM), however, reverse path forwarding (RPF) checking can fail. RPF is a validation check to determine whether a multicast packet came from the expected path, i.e., the reverse path. A RPF check can fail because the source can be an external source, i.e., a source located outside the DVMRP domain. For example, the external source can be in a PIM multicast routing domain.\nThe multicast routing information table maintained in a DVMRP router is generally limited to DVMRP-generated routing information. Typically there is no mechanism for the router to get the reverse path information of a source that is outside the DVMRP domain. Consequently, multicast traffic provided to a DVMRP multicast border router (MBR) from an external source and subsequently delivered to a DVMRP router that is not a MBR generally does not include the RPF information for the external source. As a result, the multicast traffic may be dropped.\nWhat is needed is a method to enable a non-MBR DVMRP in a DVMRP domain to perform the RPF validation check. The present invention satisfies this need and provides additional advantages."}
{"text": "Image processing is often used for photography or movie. As we know, the scanner plays an indispensable role in this field. Higher resolution is required to obtain more clear images, but the image files also occupy a larger space in the disk. An image of A4-size may be converted into a file of few megabytes when it is scanned with resolution of 600 dpi (dot per inch). Sometimes, an image with lower resolution is allowable so the size of the image file can be reduced significantly. The required resolution is various according to circumstances. For example, scanning a usual document requires a lower resolution, but scanning a beautiful photograph needs a higher resolution. Accordingly, the conventional scanner with single resolution is not suitable for all different kinds of images. Recently, the concept of a scanner with multiple resolutions has been proposed. It is conceivable that some additional elements are needed for \"switching\" among resolutions to achieve this concept.\nA driving element is usually adopted to perform the \"switching\" function. The word \"switching\" means that some optical elements is driven by the driving element to be move along with the light path according to the desired resolution of the scanning operation. It is possible but not commercially practicable because any moving action is unfavorable to the precision of the scanner. Precise locations of the optical elements are very important. IF the optical elements are moved, no one can assure that they would be moved to the right positions. A slight error occurred in moving the optical element will cause a serious error in the resulting image file. In other to improve the above-described defects, it is necessary to make efforts for manufacturing the conventional scanner with multiple resolutions.\nPlease refer to FIG. 1 which shows the operational process of a conventional scanner. The piece of paper 11 is put on the scanning window to be scanned. As shown in FIG. 1A, if a lower resolution is acceptable, the low-magnifying powered lens 121 of the lens set 12 is moved to be located within the light path so that an image received by a charge coupled device 13 can be magnified to a smaller size. On the other hand, if a higher resolution is required, the high-magnifying powered lens 122 of the lens set 12 is moved to be located within the light path and an image received by the charge coupled device 13 is magnified to a larger size, as shown in FIG. 1B. The most serious disadvantage of this method is its unsatisfied precision resulting from the movement of the lens set 12. The method also has other disadvantages. To perform this moving function, an additional space is needed for locating the driving device and moving the optical elements. It is a tendency for modem machine to have a minimum size.\nMany improved devices or methods have been disclosed. For instance, Taiwanese Patent Publication No. 281851 disclosed a method used for scanner, fax machine, copy machine or the like with optional resolutions. It teaches that both the zoom lens and the charge coupled device are movable, both of which are driven by a driving motor. In order to keep the precision of the machines, a calibrating unit is required. Besides, Taiwanese Patent Publication Nos. 286850 disclosed that the charge coupled device is moved by a ball screw. Other relevant patents, Taiwanese Patents Publication No. 232398, 235112, and 280480, teach that lenses are moved by their respective featured driving elements. The improvement of these prior publications seem nothing in comparison of their large volume, high production cost, and unsatisfied precision."}
{"text": "The present invention relates to systems and methods for improving the quality of a call over network through the optimization of packet delivery over the internet and other networks. Such systems and methods enable more efficient communications over networks than is currently available due to better communication fidelity. This results in faster and more productive calls, or any other real-time data transfer.\nCurrently, a number of platforms are available for call over network (CON) communications. These typically include audio and/or video calling capabilities that rely upon the internet (or other suitable network) in order to enable communications between two or more participants. Examples of current CON systems include Vonage, Skype, WebEx, and Facetime, for example. Each of these systems have some differentiating features, but predominantly operate via scheduling or otherwise initiating a call, and then transmitting and receiving audio and/or video material via the internet (and often other networks as well). In some of the more sophisticated CON systems, such as WebEx, participants additionally have the capability to share their computer desktop, and further, pass this functionality to other participants.\nWhile these CON systems have done a lot to enhance causal communications, the larger impact of CON systems is arguably on in relation to how businesses operate. Traditionally, in-face meetings were required to fully engage other business partners. The conference call was available, become costly when used for international calls since they operate over traditional telecommunications infrastructures. These newer CON systems have further increased how easily remote participants can communicate effectively; however, there are still a number of problems that tend to plague existing CON systems.\nFor example, proper connectivity of all users in a CON system is routinely an issue. Often one participant can have trouble joining or hearing without the other participant's knowledge. Connectivity issues are often sources of inefficiencies for call over network communications. Indeed, one of the largest problems facing CON systems is the fact that data is typically transmitted via the internet, which is a “best effort network”. Best effort network means that the data is transmitted with the best timing and quality reasonably possible. However, the data is transmitted over often torturous pathways, in sub-optimal conditions. As such, often timing and/or quality of the transmitted data are negatively impacted. Given that audio and video communications are done in real-time, these types of data transfers are very susceptible to transmission delays and/or packet loss.\nTraditional call over network systems handle this reduction in call quality and/or timing by reducing high data demanding communications. For example, in Skype, the video portion of the call may have a reduced quality, or may be halted altogether. Additionally, these existing systems simply drop the call if the timing and/or quality gets below a threshold. The thinking behind these dropped calls is that the inconvenience of not being able to communicate is less burdensome than a bad connection.\nWhile there is some merit to this business model, there are some circumstances where communication is required, even in sub-optimal network conditions. This can be especially true where a number of participants are engaging in a conference call. Here schedules are often difficult to coordinate, and as such the need to communicate via that medium, and at that time, are magnified.\nTypically, when the participants are located relatively near one another, in a location with a decent internet backbone, these measures are sufficient to ensure basic levels of communication quality. However, when participants are more remotely located (for example on different continents), these measures may simply be insufficient in order to produce decent communication.\nIndeed, for a decent voice over internet protocol (VoIP) call, total latency should be less than 400 ms, and have a packet loss of less than 2%. In contrast, the typical jitter for data between China and the United States can vary from a few hundred milliseconds to a few seconds. This is in addition to a 100-160 ms latency. Packet loss between China and the United States often varies from 0 to 100%. Clearly, the internet infrastructure is sub-optimal for communications from China to the US. Other global locations vary in their suitability for communication as well.\nAll of these drawbacks to existing CON systems require that callers repeat information more often, and reduce efficiency for all members. Moreover, in the extreme situation of a badly compromised network connection, existing CON systems are rendered inoperable. Impatient participants that have a good connection may quickly lose interest in the conversation as the pace seems unbearably slow, or as the calls are dropped. Other participants may leave the call missing much of what has been communicated. Often these participants don't want to ask for clarification too often due to respect for the numerous other participants' time. In the case of dropped calls, important communication may simply never happen.\nIt is therefore apparent that an urgent need exists for systems and methods for improving communication packet delivery over a network in order to improve the quality of a call over network. Such systems and methods provide optimization of the internet backbone, optimization of last mile delivery, and real time monitoring and mitigation technologies in order to achieve the best possible transmission outcomes given the infrastructure limitations present."}
{"text": "Electrical switch circuits are used in many applications including amplification, power, and inverter circuits. In particular, electrical switch circuits are typically used in phase switch modules. Phase switch modules are used in many applications, such as electric motor drives for industrial or automotive application to convert a direct current (hereinafter sometimes D.C.) signal to an alternating current (hereinafter sometimes A.C.) signal. Depending upon the number of electrical switch circuits in the phase switch module, a single or multiple phase switch circuit can be designed.\nA typical electrical switch circuit 100 is shown schematically in FIG. 9. The switch circuit 100 includes a first and second electrical device 102, 104, respectively. In the embodiment illustrated in FIG. 9, the devices 102, 104 are transistors. In particular, the devices 102, 104 may be, for example, Metal Oxide Semiconductor Field-Effect Transistors (hereinafter sometimes MOSFET), Insulated Gate Bipolar Transistors (hereinafter sometimes IGBT) with an antiparallel diode or the like. The devices 102, 104 illustrated in FIG. 9 are configured in a half-bridge electrical configuration. A supply terminal 106 is coupled to the drain/collector (MOSFET/IGBT) of the first electrical device 102. The source/emitter of the first device 102 is coupled to the drain/collector of the second device 104. An output terminal 108 is also coupled to the source/emitter of the first device 102 and the drain/collector of the second device 104. A reference terminal 110 (e.g., ground potential) is coupled to the source/emitter of the second device 104. Switching terminals 112, 114 are separately coupled to the gates of the devices 102, 104 through gate resistors 116, 118, respectively. Additionally, diodes 120, 122 are coupled across the drain/collector and the source/emitter of the devices 102,104, respectively. In applications using MOSFET devices, the diodes 120, 122 may be either the body diode of the MOSFET or an additional diode.\nA periodic full-wave, single phase A.C. drive signal is created on the output terminal 108 by switching the devices 102, 104. The devices 102, 104 are switched on by applying a trigger voltage to the terminals 112, 114, respectively. When the first device 102 is switched on, the second device 104 is switched off and the output terminal 108 is coupled to the supply terminal 106 through the first device 102. Conversely, when the second device 104 is switched on, the first device 102 is switched off and the output terminal 108 is coupled to the reference terminal 110 through the second device 104. Through the proper switching of devices 102 and 104 between the positive supply terminal 106 and the reference terminal 110, and by properly modulating each device's duty-cycle, the waveform created on the output terminal 108 is a positive full-wave, single phase drive signal.\nA three phase A.C. drive signal can be created in a known manner by the cooperation of three electrical switch circuits 100A, 100B, and 100C as illustrated schematically in FIG. 10. The drains/collectors of the first transistor 102A of each circuit 100A, 100B, 100C are coupled together to a positive supply terminal 106A. Similarly, the sources/emitters of the second transistor 104A of each circuit 100A, 100B, 100C are coupled together to a reference terminal 110A. A full-wave, three phase A.C. drive signal can be created by triggering each circuit 100A, 100B, 100C at alternating times. Separate full-wave, single phase drive signals are created at each output terminal 108A, 108B, 108C. Each of the three full-wave, single phase drive signals is 120° out of phase with each other. When combined, the three full-wave, single phase A.C. drive signals create a full-wave, three phase A.C. drive signal.\nA typical phase module may include one or more of the circuits illustrated in FIG. 9 and FIG. 10. In particular, some phase modules include a plurality of switch circuits 100 coupled in parallel configuration so as to increase the current capability of the phase module. The phase modules generally use a planar geometry distribution for the electrical devices of the circuits 100 in which the electrical devices are positioned on the same plane. Positioning the electrical devices on the same plane generally increases the footprint of the phase module and may require sophisticated wiring techniques such as wire bonding and laser welding. In some phase module applications, for example electric motor vehicle applications, the size and ease of manufacturability of the phase modules are important considerations."}
{"text": "The subject invention is directed toward the art of safe boxes for storage of money or other valuables to prevent theft of same and, more particularly, to safe boxes of the type mounted in the cab of a vehicle for use by the vehicle operator to deposit money collected from deliveries.\nThe invention will be described with respect to a certain preferred embodiment; however, as will become apparent, the invention can be embodied in many forms of safe boxes and enclosure structures.\nIt is known to have safe boxes of the type which can be opened only by the use of a key in the form of a special tool. After boxes of this type have been in use for some time, professional thieves may obtain the special tool required to open the box. It would be desirable to have a safe box which would be extremely difficult and time-consuming to open even by a professional thief having some special tools.\nSafe boxes or the like having a plurality of locking devices must be of special construction and the locks normally do not cooperate with one another in such a manner that sequential opening of the locking devices in a certain order is required for gaining access to the safe box. Locks which function completely independently of one another are of little benefit as a way to modify an existing single lock safe box having a lock which can be opened by some professional thieves having a special tool."}
{"text": "(1) Field of the Invention\nThis invention relates to a mold closing device of an injection molding machine. More particularly, the invention relates to a mold closing device in which the hydraulic pressure resistance in the direction of the movement of a mold closing piston is reduced so as to promote the speed of opening and closing of the injection mold.\n(2) Description of the Prior Art\nThere are several well known mold closing devices in which oil passages are formed in the mold closing piston of a mold closing cylinder and the oil passages are opened by means of valves, thereby reducing the hydraulic pressure resistance in the oil chamber on the lower pressure side when the mold closing piston is moved.\nIn almost all of them, however, either the opening or closing of valves has been done by hydraulic pressure, while it has been regarded as rather disadvantages in view of structural and economic reason for both the opening and closing of valves to be carried out by hydraulic pressure. Further, even if such an invention has been already made, it has never been put into practice."}
{"text": "Many imaging systems, such as those used in surveillance and targeting systems, include a laser or other light beam projector, such as a laser pointer, laser rangefinder, or laser designator for target designation, rangefinding, or other uses. To ensure accurate aim of the laser beam, the laser beam is typically boresighted (e.g., aligned) to an imaging sensor of the imaging system, for example, to the location of a reticle displayed in images captured by the imaging sensor. For laser designation in military weapons systems, boresighting of laser designators to imaging sensors is of significant importance for accurate delivery of ordinance. However, the alignment between the laser beam and the imaging sensor inevitably drifts over time and temperature.\nTo account for the drift in the alignment, some imaging systems are provided with a boresighting module to perform rapid, automatic boresighting periodically or prior to using the laser beam in the field. However, due to the complex optical arrangement, conventional boresighting modules are typically too complex and bulky for compact applications that fit within and/or selectively move in and out of the optical path of imaging system package. Furthermore, for a target object utilized in boresighting, conventional boresighting modules typically enclose a thermal target that burns away or otherwise deteriorates over use and must be replaced regularly. While imaging systems for surveillance and targeting often include multiples imaging sensors with different wavelength responsivity, conventional boresighting modules typically do not allow for simultaneous boresighting of multiple imaging sensors in imaging systems."}
{"text": "A trigger-type liquid ejector configured to suction a liquid from a container body and eject the liquid from a nozzle using an operation of a trigger extending downward from the nozzle is known (for example, the following Patent Document 1). In the trigger-type liquid ejector of the related art, an ejection barrel extending forward is formed on a vertical supply pipe communicating with the container body. The nozzle is attached to a tip side of the ejection barrel. A cylinder operated by the operation of the trigger is disposed below the ejection barrel. Then, when the operation of the trigger is performed, the liquid can be injected (ejected) forward from the ejection barrel via the nozzle while the liquid can be suctioned from the vertical feed tube into the cylinder."}
{"text": "Electro-hydraulic valves are used in many industrial and mobile applications. If an electro-hydraulic valve needs to be serviced or replaced, the serviced or replacement electro-hydraulic valve will need to be tuned for use with the system. Such tuning is typically done manually. However, the manual tuning of electro-hydraulic valves leads to long machine downtimes for the end user."}
{"text": "Concrete structures and other building systems often incorporate joints that accommodate movements due to thermal and/or seismic conditions. These joint systems may be positioned to extend through both interior and exterior surfaces (e.g., walls, floors, and roofs) of a building or other structure.\nIn the case of an exterior joint in an exterior wall, roof, floor, and so forth, exposed to external environmental conditions, the expansion joint system should also resist the effects of the external environment conditions. In vertical joints, such conditions will likely be in the form of rain, snow, or ice that is driven by wind. In horizontal joints, the conditions will likely be in the form of rain, standing water, snow, ice, and in some circumstances all of these at the same time. Additionally, some horizontal systems may be subjected to pedestrian and/or vehicular traffic.\nWith particular regard to bridge expansion joints, a major cause of structural deterioration of piers, columns and beams on bridges is leaking and/or deterioration of joints. Water laced with de-icing salts and atmospheric contaminants directed through expansion joints can shed directly onto critical structural elements of the bridges. Potential corrosion and subsequent spalling may occur thereby necessitating expensive reconstruction of beams, piers, columns, wing walls, and so forth.\nMoreover, expansion joint products do not fully consider the irregular nature of some expansion joints. It is common for an expansion joint to have several transition areas along the length thereof. These may be walls, parapets, columns, or other obstructions. As such, the expansion joint product follows the joint as it traverses these obstructions. In many products, this is a point of weakness, as the homogeneous nature of the product is interrupted. Methods of handling these transitions include stitching, gluing, and welding. In many situations, it is difficult or impossible to prefabricate these expansion joint transitions, as the exact details of the expansion joint and any transitions and/or dimensions may not be known at the time of manufacturing.\nAdditionally, in many products, the afore-referenced transitions present weak spots from both a water proofing aspect and a fire resistant aspect. Both expansion joints and fire resistive expansion joints typically address either water tightness aspects or the fire resistive nature, but not both. This has typically resulted in the installation of two systems for each expansion joint where both a fire rating and water resistance is required. In many cases, however, there simply is not sufficient room in the physical space occupied by the expansion joint to accommodate both a fire rated system and a waterproofing system.\nAccordingly, there exists a need for improved expansion joint systems, which can not only accommodate thermal and/or seismic movements, but also assist in alleviating and/or preventing deterioration of structural features due to environmental factors. There is a further need for such expansion joint systems that can also address fire and water resistance in one system."}
{"text": "The present application relates generally to an improved data processing apparatus and method and more specifically to mechanisms for providing a packaging for an eight-socket one-hop symmetric multiprocessing topology.\nSymmetric multiprocessing (SMP) involves a multiprocessor computer hardware and software architecture where two or more identical processors are connected to a single shared main memory. The processors have full access to all I/O devices. All processors are treated equally, with none being reserved for special purposes. Most common multiprocessor systems today use SMP architecture. In the case of multi-core processors, the SMP architecture may apply to the cores, treating them as separate processors.\nSMP systems are tightly coupled multiprocessor systems with a pool of processors running independently. The processors execute different programs and work on different data. The processors have a capability of sharing common resources (memory, I/O device, interrupt system, etc.). Processors may be interconnected using buses, crossbar switches, or on-chip mesh networks. The bottleneck in the scalability of SMP using buses or crossbar switches is the bandwidth and power consumption of the interconnect between the various processors, the memory, and the disk arrays."}
{"text": "1. Field of the Invention\nThe present invention relates to a multicolor or full-color image forming apparatus for copying machines, electrophotographic printers and the like.\n2. Related Background Art\nIn the past, electrophotographic systems, ink jet recording systems, thermal transfer recording systems and the like have been proposed as recording systems for forming a full-color image. Among them, as disclosed in U.S. Pat. No. 4,483,610, an electrophotographic system utilizing a single photosensitive member can be operated at a high speed and can use a plain paper therefore it has recently been developed more and more. In order to achieve the high speed operation, a technique of increasing the number of revolutions of the photosensitive member has been proposed. However, to further increase the speed of the operation, a system having a plurality of photosensitive members corresponding to the number of colors tends to be esteemed.\nMore specifically, as disclosed in U.S. Pat. No. 5,014,094, a recording system having four electrophotographic photosensitive members corresponding to cyan, magenta, yellow and black colors for obtaining a full-color image is suitable for the high speed operation, since it can obtain the full-color image merely by feeding a transfer sheet in a straight feeding path.\nHowever, in an image forming apparatus having a plurality of photosensitive members corresponding to the number of colors, when a mono-color image is copied from a black and white original other than a full-color original, the image forming operation including the latent image formation, development, transfer, cleaning and the like is performed only with respect to the photosensitive member associated with the black color. Thus, a surface of this photosensitive member tends to wear faster and greater than the other photosensitive members, with the result that the black color photosensitive member must be replaced more frequently than the other color photosensitive members, thereby reducing the efficiency and increasing maintenance."}
{"text": "1. Field of the Invention\nThe present invention relates to systems and methods for traffic control, and in particular to a system and method for monitoring and reserving parking spaces.\n2. Description of the Related Art\nTransportation remains a formidable problem in many urban environments. Public transportation, while often convenient, is often limited in scope and operation periods, with many starting points and destinations inadequately served. Bicycles can be used for transportation, but bicycles are inconvenient in hot, cold, or wet weather, difficult traffic conditions on roads with inadequate shoulders, or where the user wants to arrive in formal clothes. As a consequence, although their use often results in substantial traffic delays, the automobile remains the most popular form of transportation in many urban and suburban environments.\nBeyond traffic delays, there are other disadvantages with the popularity of the automobiles as a medium of transportation. Notable among these disadvantages is where to park the automobile at or near the destination of interest. Parking at many popular destinations is often difficult or impossible to find, particularly at events which draw large numbers of people at designated times (e.g. sporting events, concerts, etc.).\nA person planning a trip downtown to a theater, say, can easily use the Internet to find the address of the theater and a preferred route to reach it. What the person cannot find is a map showing parking lots near the destination, and he or she particularly cannot find information about which parking lots have or might have available parking spaces, or where there might be available parking spaces along a curb for street parking. As a result, the person must drive to a location near his or her destination and then begin a search for a parking space. He or she may be fortunate and find an inexpensive parking space quickly, or may drive for an extended period of time before settling for an undesirable parking space (distant from the destination, unsafe, and/or costly). The uncertainty in how long it will take to find a parking space also causes wasted time. A person may arrive twenty minutes early for an appointment because they felt compelled to leave their starting point early enough to assure adequate time to find a parking space, wasting time that could have been spent on other activities. Similarly, if the person does not allow enough time for any expected search for a parking space, they will arrive late at the ultimate destination. This not only results in perhaps missing a portion of the event or meeting that the person was to attend, it can also waste the time of people who were waiting for us at the destination. Wasted time can also result in wasted money.\nSearching for parking spaces also increases traffic load on many urban streets. Many cars on downtown streets are not en route to a destination, but are circling and searching for a place to park. This increases traffic congestion, energy consumption, noise, and air pollution.\nWhat is needed is an apparatus and method for providing parking information to prospective users of those spaces, and for reserving those spaces. The present invention satisfies that need."}
{"text": "More than 100 years ago, the famous physicist Ernst Abbe described a fundamental limitation of any microscope that relies on any lens or system of lenses in an imaging system to focus light or other radiation: diffraction obscures (makes fuzzy) those derails of the image that are smaller in size than approximately one-half the wavelength of the radiation. See \"Scanned-Probe Microscopes\" by H. Kumar Wickramasinghe, published in Scientific American, Vol. 261, No. 4, pp. 98-105 (October 1989). In other words, the resolution of the microscope is limited by the wavelength of thee radiation. In order to circumvent this limitation, researchers have investigator the use of, inter alia, various types of imaging probes. Scanning tunneling microscopy (hereinafter \"STM\") devices, atomic force microscopy (hereinafter \"AFM\") devices, and near-field scanning optical microscopy (hereinafter \"NSOM\") are some examples of different types of probe microscopes.\nIn STM, a metallic probe is brought sufficiently close to a conducting sample surface such that a small tunneling current is established. The magnitude of this current is extremely dependent on the tip-to-sample distance (i.e., topography of the sample surface). The tip is allowed to scan laterally across the (irregular) surface of the sample body with several angstroms separation between tip and staple in order to achieve imaging with atomic-scale resolution. The tunneling current, and hence the tip-to-sample separation, is detected and controlled by an electromechanical feedback servomechanism. In AFM, imaging is achieved in a similar manner to that of the STM except that the atomic forces (either short-range repulsive or long-range attractive) are detected instead of tunneling current. An obvious advantage to this technique is that the tip and sample do not have to be conductive, all materials exert atomic forces.\nAn NSOM device is typically comprised of an aperature locator at the tip of an elongated optical probe, the aperture having a (largest) dimension that is smaller than approximately the wavelength of the optical radiation that is being used. During device operation, the probe is positioned in close proximity to the surface of a sample body. The aperture of the probe is then allowed to scan across the surface of the sample body at distances of separation therefrom all of which distances are characterized by mutually equal forces component exerted on the probe device in the direction perpendicular to the global (overall) surface of the sample body, the scanning being detected and controlled by an electromechanical feedback servomechanism as was the case in STM and AFM.\nFor example, U.S. Pat. No. 4,604,520, describes, inter alia, a probe device having an aperture located at the tip of a cladded glass fiber that has been coated with a metallic layer. The aperture is drilled into the metallic layer at the tip of the fiber at a location that is coaxed with the fiber. The (immediate) neighborhood of the tip is composed of a section of solid glass fiber that has obliquely sloping (truncated conical) sidewalls, whereby the sidewalls do not form a cylinder of any kind. Therefore, as the probe device laterally scans a rough surface, the calculations required to determine the desired information on the actual contours (actual profile) of the surface of the sample body require prior detailed knowledge of the slanting contours of the sidewalls of the probe, and these calculations typically do not yield accurate metrological determinations of the desired profile of the contours of the surface of the sample body, especially at locations of the surface of the sample body where sudden jumps (vertical jumps) thereof are located. In addition, fabrication of the probe device is complex and expensive, especially because of the need for drilling the aperture coaxially with the fiber.\nAnother example involves the fabrication of nanometric tip diameter fiber probes for photon tunneling microscopes (\"PSTM\") by selective chemical etching of the GeO.sub.2 -doped cores of optical fibers. See \"Reproducible Fabrication Technique of Nanometric Tip Diameter Fiber Probe for Photon Scanning Tunneling Microscope\", Togar Pangaribuan, et.al., Japan Journal Applied Physics, Vol. 31 (1992), pp. L 1302-L 1304. By selectively etching the GeO.sub.2 doped regions of the fiber, a tapered tip having the shape of a small cone can be formed on the endface of the optical fiber. The cone angle of the fiber probe tip is controlled by varying the doping ratio of the fiber core and the composition of the etching solution. A fiber probe with a cone angle of 20.degree. and tip diameter of less than 10 nm was fabricated. Only probes having conical-shaped endfaces can be made with this technique, so that the sidewalls do not form a cylinder of any kind. The scanning range of such a probe is undesirably limited owing to the relatively large width (diameter) of the endface on which the relatively short-width conical tip is centered, coupled with the fact that, during scanning, the probe is rastered from side-to-side in an arc: a desired large length of scan is attempted, the corners of the probe's endface undesirably will make contact with the sample surface. In addition, the conical shape of the tip undesirably limits the accuracy of measurements wherever the surface being probed has a sudden jump."}
{"text": "Variations in the electrical conductance of the intact skin of human beings has been recognized for many years. Particularly high conductance sites have been associated with acupuncture points as well as with trigger points in myofascial pain syndromes. However, the realization that this change in electrical skin conductance was related to the presence of subcutaneous nerves was only made recently by the inventor, and represents a unique finding. Although the measurement of transcutaneous: electrical conductivity has been used for other purposes, it has never before been shown to reveal the position of neural tissue.\nThe present non-invasive-peripheral nerve mapping device and method may be used to indicate the site of nerve injury or nerve transection. Sudden loss of transcutaneous electrical hyperconductivity along the course of a peripheral nerve is indicative of interruption of the nerve fiber whereas a focal increase in the transcutaneous electrical hyperconductivity along the course of a peripheral nerve is evidence of a site of nerve injury without actual transection."}
{"text": "1. Field of the Invention\nThe present invention relates to a semiconductor structure, and more particularly, to an interconnection structure.\n2. Description of the Related Art\nAfter the semiconductor fabricating process has developed into the deep micron field, some originally insignificant issues have become obvious problems due to the reduced size. In order to save space for further reducing the chip size, the electrostatic discharge circuit is commonly disposed below the pad.\nDuring a wire bonding process using a thermocompression method, since too much pressure is applied on the pad, the dielectric layer sometimes cracks or even the pad would come off, leading to poor yield and low reliability of the package production. The cause of such problem is insufficient adherence between the dielectric layer and the metal layer and the stress from wire bonding process applied on the dielectric layer.\nIn order to prevent the semiconductor element below the pad from such impact and to ensure the successful bonding, it is quite often that the pad uses two metal layers, so the rest of the lower metal layers can be used for circuit wiring. For the fabricating process with fewer metal layers, the number of the metal layers used for circuit wiring is insufficient."}
{"text": "1. Field of the Invention\nThe present invention relates to organ-specific molecular blood fingerprints and methods for using the same in identifying and/or monitoring drug side effects.\n2. Description of the Related Art\nSide effects, particularly adverse side effects, are monitored and tested throughout the developmental path of all drugs. A variety of in vitro and in vivo toxicity assays are available in the art for testing on and off-target effects of drugs during development including metabolic effects on liver P450 enzymes. Clearly, however, as is evidenced by the recent controversy surrounding COX-2 inhibitors, off-target side effects can be subtle and difficult to detect. Therefore, monitoring and identifying drug side effects remains an important issue with regard to drug safety of developing drugs and approved drugs already on the market. The COX-2 story highlights the need in the art for improved methods to detect sometimes subtle adverse off-target effects of drugs.\nThe present invention provides methods that satisfy this and other needs."}
{"text": "The invention pertains to a device for projection-copying of a mask onto a workpiece, in particular a semiconductor substrate for the production of integrated circuits, in which case the pattern of the mask is projected via a projection lens onto a photosensitive layer of the workpiece after in order to align the mask and workpiece relative to one another alignment patterns of the mask and adjustment areas on the workpiece have been projected onto each other by means of an adjustment light with a bandwidth of at least 3 nm through the projection lens.\nProjection lenses for the lithographic production of integrated circuits are charcterized by a large picture field with a diameter which typically lies between 10 and 30 mn and a large numerical aperture when there is a great resolution capability which is limited with respect to diffraction. Because of the necessity of projecting different pictures one on top of the other in such a way that they precisely cover one another, the distortion in the entire field must not exceed 0.2 .mu.m and the picture field must be completely level, i.e., any convexity must not exceed 0.3 .mu.m.\nAt the present state of the art, lenses which meet such stringent requirements can be corrected up to the diffraction limit only for a very narrow wavelength range, which is to be understood as a bandwidth of a few nm.\nThis narrow correction area is selected in such a way that within it the sensitivity of the photosensitive resist is as high as possible and, on the other hand, an appropriate source of illumination is available. Typical correction ranges are 406 nm.+-.4 nm or 436 nm.+-.4 nm, corresponding to the most intense spectral lines of mercury discharge lamps.\nAs mentioned above, it is vital for the projection of the pattern of the mask to take place not only with good picture quality, but also with complete precision of positioning. The precision with respect to the lateral coordinates (X, Y, .theta.) is necessary in this case in order for the successive patterns to be correctly allocated, but in addition it is also necessary to precisely focus an entire picture plane since the depth of focus of the above-described lenses is very slight.\nThe alignment of the mask and workpiece is preferably carried out through the lens itself, and in this process the adjustment areas for lateral adjustment are defined by marks of the most widely varying structure. Per se, the unaltered reflective surface of the workpiece itself is sufficient for focusing.\nFocusing as such does not necessarily require that the workpiece be observed through the lens. For instance, it is a known process to determine the distance between the lens and the workpiece by means of capacitive sensors with the aid of the flow resistance which the annular gap between the lens and the workpiece presents to a discharging gas, or from the eigen-frequency of an air resonance section between the workpiece and the lens. A prerequisite in this case, however, is that there be a very short distance between the lens and the workpiece, but such a short distance is precisely what is avoided through the use of projection illumination procedures, in contrast to the obsolete contact procedures. When optical focusing is carried out without the use of the projection lens, a tightly packed beam of light is obliquely directed onto the center of the field to be illuminated and the point where the reflected light beam strikes is observed; the location of this point is a measure of the location of the illumination field itself. The light used is frequently laser light (a HeNe laser). A disadvantage in this case is the fact that it is impossible to distinguish whether the change in the position of the light point on the receiver is caused by a change in the angle of incidence or the Z position. In practice, it is assumed that the angle of incidence is constant, i.e., that the reflective wafer surface is always vertical to the optical axis. This condition is never rigorously met; sometimes there are even large deviations from the ideal position. When laser light is used, another problem arises due to diffraction effects at the wafer surface (speckle), in particular if the surface has already been structured (higher manufacturing stages).\nSince, when adjustment is done through the projection lens, the system-related disadvantages of the above-described procedures do not arise, there is great interest in solving the sub-problems which still exist with this type of adjustment. When adjustment is done through the projection lens, problems arise when the adjustment light is not identical to the illumination light for which the projection lens is corrected. Initially we think here of the case where the wavelength range of the adjustment light lies outside of the range of the spectral sensitivity of the photosensitive resist to avoid having the marks on the workpiece be destroyed by the adjustment process. The differences with respect to focal length and magnification which the lens shows depending on the type of light used can be compensated for by bending the adjustment beam through a pair of mirrors or by lengthening the beam by means of intermediate glasses or by shifting the location of the protector for adjustment. Since tests can readily determine the extent to which the focal length and magnification of the lens differ at the corresponding illumination wavelength and adjustment wavelength, overall the consequences of this difference can be easily handled; with respect to the position where the device is optimally aligned with the adjustment light, prior to illumination being carried out a shift is simply made which takes the differing behavior of the lens in the two cases into account.\nSince, with an illumination wavelength which deviates from the adjustment wavelength, the picture defects of the lens can be corrected only when the bandwidth of the adjustment light is relatively small, it is assumed that the adjustment light used should, in principle, be narrowband. The frequenlty-made suggestion that laser light be used for adjustment purposes has not proven out in practice since, due to the coherence of this light, diffraction effects (speckle) occur which distort the measurement results. Generally mercury discharge lamp light, the natural line width of which is approximately 3 nm, is thus used for adjustment.\nSurprisingly enough, it has been found that, when the adjustment is carried out with mercury light through the projection lens in the way described in the introduction, a nonsystematic error arises, the cause of which was found to be the fact that the reflection capacity of the workpiece is dependent on the processing stage of the workpiece and, in addition, that this capacity varies on the surface of the workpiece. Not only does the color of the reflected adjustment light deviate slightly from the color of the input adjustment light, but differences in the color of the reflected adjustment light appear from workpiece to workpiece and from mark to mark on the same workpiece. If it is assumed that, in addition to the spectral line itself, the adjacent area of the radiation background in a total width of, for example, 10 nm is also typically passed by the narrowband interference filter in front of the adjustment light sources, it is still amazing that the differences of 1-2 nm which arise overall in the wavelength of the reflected adjustment light still have an effect on the precision of the adjustment. Due to the heavy frequency dependency of the lens used, on the one hand, and the extreme demands imposed on adjustment precision, on the other, however, this is indeed the case.\nThe invention thus is based on the recognition that it is not sufficient, as was previously the assumption, to take into account the differences which arise, regarding the focal length and magnification of the lens, in the illumination wavelength on the one hand and, on the other, and adjustment wavelength which is assumed to be constant, but rather it is necessary to take into account the effect of the difference in the reflection behavior of the workpieces, which difference cannot be known in advance, and which leads to a change in the spectral composition of the input adjustment light."}
{"text": "Conventionally, movable molds are employed in continuous casting of thin metallic plates. Two types of such molds are known: one for drawing a casting in a generally horizontal direction, and the other for drawing a casting in a vertical direction. The former, as FIG. 29 illustrates, is a belt-type movable mold 103 positioned under a melt receiver 101 and endlessly trained around a pair of drive rollers 102A, 102B, front and rear, in a longitudinally inclined condition. Molten metal within the receiver 101 is drawn in the form of a thin plate-like casting sequentially as it is cooled and solidified on the surface of the belt-type mold 103. A movable mold of the latter type, as FIG. 30 shows, is a twin-roll type mold 113 wherein molten metal is drawn through a clearance defined between a pair of rolls 112A, 112B disposed under a melt receiver 111 and is then solidified in the form of a thin plate-like casting. Aforesaid melt receivers 101 and 111 are of a rectangular configuration, each including a pair of side walls 101a or 111a, as the case may be, disposed parallel to the path of casting withdrawal, and a rear wall 101b or front and rear walls 111b, 111b, as the case may be, disposed in rectangular relation to the withdrawal path.\nSuch conventional melt receiver is held stationary at a specified position and, therefore, it has a drawback that a casting drawn is often liable to break-out. The cause of such break-out may be explained as follows. A wall shell is formed at a corner portion defined between the movable mold and either one of the side walls or the front or rear wall, and it tends to grow. A grown wall shell adheres to a casting shell formed on the movable mold through solidification of molten metal and tends to pull the casting shell away from the direction of drawing, with the result that the casting shell is broken at a relatively thin portion thereof. Again, such wall shell adheres to the casting shell until such breaking occurs. This process of adhering and breaking is repeated again and again, with the result that the casting shell is not constant in thickness, being thus liable to break-out at its relatively thin and weak portions."}
{"text": "1. Field of the Invention\nThe present invention generally relates to a thin film transistor (TFT) array substrate, in particular, to a TFT array substrate with capacitance compensating lines.\n2. Description of Related Art\nIn recent years, the development of liquid crystal display television (LCD-TV) has been developed toward high resolution and large size. In a large-sized LCD panel, signal distortion may be caused due to overall RC delay of the LCD panel. Accordingly, the driver IC in a large-sized LCD panel has to be designed with strong driving capability in order to maintain the display uniformity of the LCD panel at a specific standard. On the other hand, in a high-resolution display panel having a driving frequency of 120 Hz (Full HD; 1920×1080), overheat issue on the driver IC may be caused because of the high polarity conversion frequency when a general dot inversion driving method is adopted to improve the image quality of the display panel. In this case, overheat issue on the driver IC can be significantly reduced if the polarity conversion frequency of the driver IC can be reduced. Accordingly, the problems of overheat issue on the driver IC and RC delay can be resolved when a column inversion driving method can be adopted for driving the display panel.\nHowever, there are still some other problems in image presentation to be resolved even if the column inversion driving method is adopted. For example, bright/dark lines and serious cross-talk may be produced when dynamic images are displayed. To resolve these problems, an alternative type pixel layout is disclosed in U.S. Pat. No. 6,822,718, when the display panel having the above-mentioned alternative type pixel layout is driven through column inversion driving method, the display quality thereof is similar with the display quality provided by display panel that is driven through dot inversion driving method. Since the problem of overheat issue on driver IC will not be produced by the column inversion driving method, the problems of bright/dark lines, cross-talk, and overheat issue on driver IC can all be effectively resolved by the U.S. Pat. No. 6,822,718. However, it should be noted that in the pixel layout provided by the U.S. Pat. No. 6,822,718, TFTs are disposed alternatively and which reduces the regularity of the photomask pattern used for fabricating the TFTs. Accordingly, mis-aligment may occur and black matrix having sufficient area is required to prevent light-leakage. As a result, both the aperture ratio of the display area and the size of the display area are reduced."}
{"text": "Technical Field\nThe present disclosure relates to a multi-protocol determining method. More particularly, the present disclosure relates to a multi-protocol determining method based on a control area network (CAN) bus.\nDescription of Related Art\nThere are many types of commercial vehicles, such as light or large buses, light or heavy truck, and estate cars. Such vehicles are produced in USA, Japan or Europe, and use ISO protocol or SAE protocol.\nHowever, new vehicles produced in USA are regulated to use the CAN bus interface from 2008. The vehicles have to provide an On Board Diagnostic II (OBD-II) connector, and the heavy duty vehicles has to be communicated based on SAE-J1939 while other vehicles are communicated based on ISO-15031/SAE-J1979. Moreover, new protocol, WWH-OBDII, has been developed, and some new vehicles can be communicated based on ISO-27145. Hence, how to communicate with vehicles having different protocols is a problem needed to be solved.\nSome practitioners develop a gateway which can be connected between two CAN buses, i.e., the first CAN bus and the second CAN bus, having different protocols. The gateway can receive a message from the first CAN bus with one protocol and translate the message based on the other protocol. Subsequently, the gateway can send the translated message to the second CAN bus. Since the protocols used in both the first CAN bus and the second CAN bus have to be known before translation, the usage thereof is limited.\nTherefore, how to identify protocols of a communicated node of the vehicle or of a CAN bus becomes a pursuit target."}
{"text": "Present processes are known for the manufacturing of both fine-thickness, i.e. fine and normal thickness particle board of fine and/or normal sized chips. A first known type of process relates to a continuous-flow production of particle board using calenders and the like, for example \"Mendes\" process. The second known type of process relates to a continuous production of particle board using a conveyer-belt, which moves first from under a filling machine, where the mat is made and then carries the mat under the press, where the final pressing operation occurs. This type of operation which is somewhat similar to the type of the present invention may be carried out in two different manners.\nA first and better-known process essentially employs a movable filling machine and a fixed press. A conveyor-belt, either formed of a plastic-covered canvas and/or similar material, or formed of a steel-belt, moves between the filling machine and the fixed press and carries the mat from the former to the latter. To obtain the finished particle board, the mat is compressed while it lies on the belt at the press stage.\nIf a steel belt is used in this process, a number of difficulties are encountered due to the considerable expansion of the belt during the pressing operation. Furthermore, another disadvantage with the use of a metallic belt is the fact that it becomes impossible for the superheated vapor, generated within the pressed material, to be released satisfactorily.\nOn the other hand, if a fabric or similar belt is employed to eliminate these defects which are encountered with the use of a metallic belt, the produced particle board is often insufficiently flat, the thickness is inconsistent and the surface finish is not smooth and even. When this occurs, it is necessary to further finish the particle board, which results in increased manufacturing labor and costs. Also, this type of belt is considerably less durable than metallic belts, thereby entailing high maintenance costs.\nA second process employs a belt system shuttling between the filling machine and the press, the belt depositing the mat directly on the press plate so as to obtain a good tolerance and surface finish quality of the final product and to eliminate subsequent surface finishing. Furthermore, this system obtains also a superior quality product. However, the system is extremely complicated and it contains the defects mentioned above and, futhermore, cannot be used to obtain fine thicknesses. In this case great difficulties are encountered in depositing the material in the plane of the press for producing thin products.\nProcesses with a movable filling machine and a movable press, while suggested, have not been in use in view of the difficulties in shifting these machines."}
{"text": "It is known in the art to compress a database including data to minimize storage requirements for storing the data and reduce transmission times for transmitting the data. In prior approaches, the entire database is compressed and decompressed or extracted for manipulation/query of the data in the database. For example, prior approaches are directed to reducing the search time required for searching over a large database using methods such as binary searches or b-trees both of which require that the data in the database can be read randomly. In order to support random reading from and writing to a compressed database, the entire database must be decompressed.\nThere is a need in the art for a method of enabling update of data in a database having a compressed data structure without requiring decompression of the entire database. That is, the database remains compressed and occupies a smaller storage space thereby requiring less memory and less transmission time to transfer the database update contents.\nFor example, handheld or embedded devices are constrained by limited processing power and limited storage or memory in order to increase the device's battery life. A method of updating a compressed database would enable a larger amount of data to be stored on the device and would decrease update time. However, prior approaches have always decompressed the entirety of the data prior to update of the data on the device thereby eliminating any advantage gained from database compression."}
{"text": "1. Field of the Invention\nThis invention relates to a cover assembly for an electronic device, specifically including a laptop computer, having a base removably mounted in an operative, protective position on the electronic device and movable therewith between open and closed positions. A retaining assembly is disposed at least partially on the interior of the electronic device and is structured to facilitate selective and complete closure thereof. The base is structured to provide access to the connecting ports and other operative components of the electronic device, whether it is in the open or closed positions.\n2. Description of the Related Art\nVarious electronic devices, including laptop or notebook computers are currently being used in commercial, domestic and educational environments. In particular, the younger generation has made ubiquitous use of small, lightweight portable computers in a school or educational environment, wherein such portable computers are carried from class to class. The portability of the laptop computer also lends itself to being carried or utilized during travel on both private and commercial modes of transportation including airlines, automobiles and the like.\nTherefore, laptop computers are designed for substantial portability thereby facilitating an individual's access in almost an unlimited number of locations. As such, one problem commonly associated with the transport and use of a portable computer is the requirement for at least some minimal protection against situations where the computer may be dropped or otherwise be subjected to other traumatic forces. In order to overcome such potential problems, there have been developed a variety of different “carrying case” structures. These known structures are similar to luggage type products, inclusive of shoulder straps, handles, etc.\nThese known carrying cases are available in a wide variety of styles and sizes and are used not only to protect the computer against damage but also to facilitate the transport of the computer in a convenient and comfortable manner. Also, many of the known carrying case devices are structured to protect the portable computer against adverse weather conditions, as well as damage from being dropped, etc. However, many of the known structures of this type are bulky, somewhat heavy and commonly involve relatively complicated enclosure or locking assemblies. As such, known devices of this type mandate that the computer be removed from the interior of the device, placed on an appropriate support structure then disposed in an open orientation for use.\nAccordingly, there is a long existing need in the area of cover assemblies for portable computers, which provide protection from both adverse weather conditions as well as traumatic impact. Such a proposed cover assembly should be lightweight, easily mounted in an operative, protective position on the laptop computer in an efficient manner. Further, such a proposed cover assembly should be capable of maintaining the preferred operative position in protective, covering relation to the portable computer or other electronic device during its use, storage or transport. More specifically, the operative position of the proposed cover assembly is maintained when the computer, etc. is in an open position, such as when it is being used, as well as a closed position, such as when the computer is being stored or carried. The ability to maintain the proposed cover assembly in its operative position thereby eliminates the necessity of removing it in order to use the portable computer as is common place with most protective, carrying case type of structures, now available.\nTherefore, a proposed cover assembly of the present invention should be formed of a lightweight material having capabilities which facilitate the protection of the computer from adverse weather conditions, as well traumatic force or impact. Moreover, the various structural and operative features of such a proposed cover assembly should enable the portable computer or like device to be used in an intended, conventional manner. Such conventional use would include it being opened and fully closed, while certain retaining structures associated with the proposed cover assembly remain in place. Thereby the proposed cover assembly would be securely but removably disposed in its intended operative position without interfering with the orientation or use of the portable computer or other electronic device."}
{"text": "A user device may connect with a network element to allow a user to perform a task via the user device (e.g., browse the web, send and/or receive e-mail, etc.). Data flow, associated with performing tasks via the user device, may be cumbersome and redundant when the user device relocates to another geographic location and changes network element connections."}
{"text": "1. Field of the Invention\nThe embodiments described herein relate generally to inspection of containers and, more particularly, to producing images of at least one object of interest in a container to facilitate detecting objects within the container.\n2. Description of Prior/Related Art\nAt least some known inspection systems construct an image of a container and analyze the image to detect explosives, drugs, weapons, and/or other contraband objects within the container. Some known inspection systems, including some known carry-on and checked baggage inspection systems, use a computed tomography (CT) scanner to produce images of the interior of the container. In some known CT systems, an X-ray fan beam source and a detector array are disposed opposite each other in a gantry. The gantry is rotated around the container such that the angle at which the X-ray fan beam intersects the container constantly changes. The detector array acquires a “view,” or group of X-ray attenuation measurements, at each of multiple gantry angles. A “scan” of the object includes a set of views made at different gantry angles, or view angles, during one revolution of the X-ray source and detector about the container.\nAt least some known CT systems use three-dimensional (3-D) volumetric scanners to acquire CT scan data representing an entire target volume. One example of a 3-D volumetric CT scanner is a helical CT scanner, in which the container is continually moved substantially parallel to the axis of gantry rotation while the X-ray attenuation data is being acquired, such that the path of the X-ray source defines a helix with respect to the container. Three-dimensional representations of the entire volume scanned by a 3-D volumetric scanner can be reconstructed using well known tomographic reconstruction algorithms, for example, direct Fourier or filtered back-projection methods.\nAt least some known inspection systems also include a prescanner to produce a two-dimensional scan projection (SP) image of the container for presentation to an operator for inspection. Some known prescanners operate by moving the container under a fan beam of X-rays from a stationary X-ray source. X-ray intensities, after being attenuated by the container and the objects within it, are measured by a stationary array of detectors. Such prescanners may use a separate stationary X-ray source and detector array, or they may employ a primary 3-D volumetric CT scanner in a stationary prescanner mode to produce two-dimensional SP images before conducting a helical scan. Inspection system operators typically receive extensive training, and/or have accumulated extensive experience, in recognizing certain types of objects in a two-dimensional SP image. As a result, for at least some inspection system operators, inspection of an SP image, rather than of a full three-dimensional representation from the CT scan, increases a speed and accuracy of the identification of certain types of objects.\nThe use of a dedicated prescanner also may provide other benefits. For example, at least some known prescanners provide dual energy scanning of the container. Collecting data for a low-energy scan and a high-energy scan allows such inspection systems to reconstruct, for example, a density image and/or atomic number image of the contents of the container to facilitate identification of objects and materials in the container.\nUnfortunately, the use of a prescanner also may increase a cost and complexity of the inspection system, for example by requiring a separate stationary X-ray source and detector array, or by requiring additional hardware and software in the 3-D volumetric CT scanner to support a stationary prescanner mode. The time required to conduct the prescan also adds to the overall time required for scanning a container. Some known inspection systems avoid these potential drawbacks by forgoing a prescanner, instead re-projecting the 3-D volumetric CT scan data into two-dimensional images similar to SP images from a prescanner. However, in at least some such systems, a resolution of the re-projected two-dimensional SP-type images is less than a resolution of the SP images produced by a prescanner, and/or less than a resolution of the original three-dimensional CT representation. This reduced resolution may cause objects of interest identifiable from the three-dimensional CT scan data, such as thin wires, not to be visible in the re-projected two-dimensional image.\nFor example, FIG. 3 is a right perspective view of certain representative objects of interest that may lie within a container 18. In FIG. 3, container 18 and any non-numbered objects within container 18 are rendered transparent for ease of viewing. Container 18 might include objects of interest such as, for example but not by way of limitation, a thin wire 60 between two otherwise non-descript objects 61 and 62, a plastic explosive material 64, and/or overlapping objects 66 and 68.\nFIG. 4 shows a prior art two-dimensional SP-type image that might be obtained using a suitable method for re-projecting an SP-type image from data obtained from a 3-D volumetric scan of container 18 containing representative objects of interest 60, 64, 66 and 68, as shown in FIG. 3. Alternatively, an image such as that illustrated in FIG. 4 might be obtained using a suitable method for re-projecting an SP-type image from a three-dimensional representation of container 18 that was in turn generated from 3-D volumetric scan data from a scan of container 18.\nWith reference to FIGS. 3 and 4, otherwise non-descript objects 61 and 62 appear in the SP-type image as disconnected areas 71 and 72, with no indication of thin wire 60, although thin wire 60 would have been detected in the inspection of the three-dimensional image representation. Similarly, plastic explosive 64 appears as a non-descript region 74 in the SP-type image, although in certain embodiments a multiple-energy scan by inspection system 10 would have detected a material comprising plastic explosive 64 as a material of interest. In addition, overlapping objects 66 and 68 appear as a single unified region 76 in the SP-type image, although the presence of separate objects 66 and 68 would have been detected in the inspection of the three-dimensional image representation."}
{"text": "1. Field of the Invention\nThis invention relates to microprocessors and, more particularly, to branch prediction mechanisms employed by microprocessors.\n2. Description of the Related Art\nMany high performance pipelined microprocessors maintain their performance by keeping the various pipelines as full as possible. However, certain control transfer instructions may cause various stages within a pipeline to be flushed and the instructions in those stages to be canceled or discarded. One such instruction is a branch instruction. A branch instruction is an instruction which causes subsequent instructions to be fetched from a target address identifying an instruction stream beginning at an arbitrary location in memory. Unconditional branch instructions always branch to the target address, while conditional branch instructions may select either the next sequential address or the target address based on a condition such as the outcome of a prior instruction. Thus, when instructions are prefetched into a processor's instruction-processing pipeline sequentially, and a conditional branch is taken, the contents of early stages of the pipeline may contain instructions that should not be executed but rather should be flushed.\nAccordingly, in conjunction with the pre-fetching of instructions, it is generally beneficial to predict whether a conditional branch instruction, when executed, will be taken or not. For this purpose, many pipelined microprocessors employ a branch prediction mechanism.\nThe branch prediction mechanism may indicate a predicted direction (taken or not-taken) for a branch instruction, allowing subsequent instruction fetching to continue within the predicted instruction stream indicated by the branch prediction. Instructions from the predicted instruction stream may be speculatively and/or executed prior to execution of the branch instruction, and are placed into the instruction-processing pipeline prior to execution of the branch instruction. If the predicted instruction stream is correct (which is generally well in excess of 90% for many current designs), then the efficiency of instruction execution may be advantageously improved. However, if the predicted instruction stream is incorrect (i.e. one or more branch instructions are predicted incorrectly), then the instructions from the incorrectly predicted instruction stream are discarded from the instruction-processing pipeline and the efficiency of instruction execution may be degraded.\nTo be effective, the branch prediction mechanism must be highly accurate such that the predicted instruction stream is correct as often as possible. Frequently, a history of prior executions of a branch or branches is used to form a more accurate prediction for a particular branch. Such a branch prediction history typically requires maintaining data corresponding to the branch instruction in a storage.\nMany branch prediction mechanisms have been implemented that yield relatively good results. One such branch prediction mechanism is the GShare branch prediction mechanism. The GShare mechanism combines a global branch history and several bits of the program counter (PC) of a current branch instruction (e.g., fetch address of the current branch instruction) to form an index into a predictor table. Each entry in the predictor table typically stores a counter value that may be indicative of whether a branch should be taken or not. Although the GShare mechanism may be an effective prediction mechanism, aliasing of index values is a concern. For example, the GShare mechanism may produce the same index value and therefore an erroneous. prediction, when provided with different input combinations of the PC and global branch history that correspond to different branches."}
{"text": "Assay test kits are currently available for testing for a wide variety of medical and environmental conditions or compounds, such as a hormone, a metabolite, a toxin, or a pathogen-derived antigen. FIG. 1 shows a typical lateral flow test strip 10 that includes a sample receiving zone 12, a labeling zone 14, a detection zone 15, and an absorbent zone 20 on a common substrate 22. These zones 12-20 typically are made of a material (e.g., chemically-treated nitrocellulose) that allows fluid to flow from the sample receiving zone 12 to the absorbent zone 22 by capillary action. The detection zone 15 includes a test region 16 for detecting the presence of a target analyte in a fluid sample and a control region 18 for indicating the completion of an assay test.\nFIGS. 2A and 2B show an assay performed by an exemplary implementation of the test strip 10. A fluid sample 24 (e.g., blood, urine, or saliva) is applied to the sample receiving zone 12. In the example shown in FIGS. 2A and 2B, the fluid sample 24 includes a target analyte 26 (i.e., a molecule or compound that can be assayed by the test strip 10). Capillary action draws the liquid sample 24 downstream into the labeling zone 14, which contains a substance 28 for indirect labeling of the target analyte 26. In the illustrated example, the labeling substance 28 consists of an immunoglobulin 30 with a detectable particle 32 (e.g., a reflective colloidal gold or silver particle). The immunoglobulin 30 specifically binds the target analyte 26 to form a labeled target analyte complex. In some other implementations, the labeling substance 28 is a non-immunoglobulin labeled compound that specifically binds the target analyte 26 to form a labeled target analyte complex.\nThe labeled target analyte complexes, along with excess quantities of the labeling substance, are carried along the lateral flow path into the test region 16, which contains immobilized compounds 34 that are capable of specifically binding the target analyte 26. In the illustrated example, the immobilized compounds 34 are immunoglobulins that specifically bind the labeled target analyte complexes and thereby retain the labeled target analyte complexes in the test region 16. The presence of the labeled analyte in the sample typically is evidenced by a visually detectable coloring of the test region 16 that appears as a result of the accumulation of the labeling substance in the test region 16.\nThe control region 18 typically is designed to indicate that an assay has been performed to completion. Compounds 35 in the control region 18 bind and retain the labeling substance 28. The labeling substance 28 typically becomes visible in the control region 18 after a sufficient quantity of the labeling substance 28 has accumulated. When the target analyte 26 is not present in the sample, the test region 16 will not be colored, whereas the control region 18 will be colored to indicate that assay has been performed. The absorbent zone 20 captures excess quantities of the fluid sample 24.\nIn the non-competitive-type of lateral flow assay test strip designs shown in FIGS. 2A and 2B, an increase in the concentration of the analyte in the sample results in an increase in the concentration of labels in the test region. Conversely, in competitive-type of lateral flow assay test strip designs, an increase in the concentration of the analyte in the fluid sample results in a decrease in the concentration of labels in the test region.\nAlthough visual inspection of lateral flow assay devices of the type described above are able to provide qualitative assay results, such a method of reading these types of devices is unable to provide quantitative assay measurements and therefore is prone to interpretation errors. Automated and semi-automated lateral flow assay readers have been developed in an effort to overcome this deficiency.\nIn one approach, a portable lateral flow assay reader performs assays on bodily fluids to detect the presence of certain hormones, glucose, or other bodily fluids of interest. Membrane test strips containing a fluid sample are inserted directly into a receiving port of a reader. The receiving port is shielded to improve sensitivity and reduce the entry of stray or ambient light into the reader. The reader includes a light source and one or more sensors that detect the intensity of light reflected from the detection zones of the test strips that are inserted into the receiving port.\nIn another approach, a reader detects an intensity of a detection signal arising in one or more measurement zones in a detection zone of a lateral flow assay test strip as a result of the presence of an immobilized labeled target analyte complex. The reader generates a baseline of signal intensity from the measurement zones by interpolating between values of the detection signal outside of the measurement zones and inside of the detection zone. The reader quantifies a value of signal intensity representative of the immobilized labeled target analyte complex with respect to the baseline. In this process, the reader locates a beginning boundary and an ending boundary for the one or more measurement zones on the test strip, allowing an automatic or semi-automatic analytical instrument, or a human reader, to determine certain results of the lateral flow assay. The signals from the measurement zones are quantified or compared with respect to the baseline. Quantified values corresponding to the respective concentration of compounds in different measurement zones may then be compared with one another to detect the presence of antigens in the sample.\nThe measurements that are made by the above-described lateral flow assay readers are based on signals from regions of the test strips that typically are significantly larger than the regions of interest. As a result, these measurements tend to have high noise levels and, consequently, these measurements may yield inaccurate or incorrect results when low concentrations of analytes are involved."}
{"text": "1. Technical Field\nThe present invention relates to a cartridge which accommodates a printing material therein, a printing material supply system which includes a cartridge and a printing apparatus, and printing apparatus to which the cartridge is detachably mounted.\n2. Related Art\nA printer which is an example of a printing apparatus has a cartridge that accommodates ink therein and is detachably mounted thereto, and performs printing by supplying the ink from the cartridge to a printing head. Recently, a cartridge in which a circuit board that stores information regarding ink (for example, the color of ink) is mounted is used (for example, JP-A-2003-341100, JP-A-2007-230249, Specification of U.S. Pat. No. 6,203,147, JP-T-2002-505212, and JP-T-2002-513341). The circuit board mounted in the cartridge includes a storage device that stores the information regarding the ink and a terminal group (also referred to as a “cartridge side terminal group”) for electrical connection to a printer. As the cartridge is mounted to a cartridge mounting unit of the printer, the cartridge side terminal group comes into contact with a terminal group (also referred to as an “apparatus side terminal group”) provided in the cartridge mounting unit. Accordingly, each terminal of the cartridge side terminal group is electrically connected to the corresponding terminal of the apparatus side terminal group, such that signal reception and transmission between the printer and the cartridge is performed.\nHowever, in the techniques of JP-A-2003-341100 and JP-A-2007-230249, due to manufacturing errors in the cartridge mounting unit or the cartridge, there is a possibility that the cartridge side terminal group and the apparatus side terminal group may not come into contact with each other so as to be electrically connected.\nIn the technique of Specification of U.S. Pat. No. 6,203,147, a connector provided with the apparatus side terminal group is movably provided in the cartridge mounting unit. Specifically, the connector is configured so as to be movable in a direction (the X axis direction and the Y axis direction in FIG. 7 of Specification of U.S. Pat. No. 6,203,147), orthogonal to an insertion direction (the Z axis direction in FIG. 7 of Specification of U.S. Pat. No. 6,203,147), of the cartridge. In addition, in the technique of Specification of U.S. Pat. No. 6,203,147, a recessed portion is provided in the front surface of the cartridge, and the cartridge side terminal group is arranged on the side surface in the recessed portion along the insertion direction. In this technique, since the connector is configured to be movable, the connector is moved according to the position of the cartridge when the cartridge is mounted to the cartridge mounting unit. Accordingly, there is an attempt to achieve proper connection between the cartridge side terminal group and the apparatus side terminal group.\nIn the technique described above, positioning in the X axis direction and the Y axis direction of the connector is performed by four side surfaces of the recessed portion in which the cartridge side terminal group is arranged. Therefore, errors in the shape of the recessed portion occur due to design errors in the cartridge, and thus there is a possibility that contact between the cartridge side terminal group and the apparatus side terminal group may not be properly achieved.\nIn addition, in the technique described above, there is a possibility that the apparatus side terminal group may rub against the side surface of the recessed portion or the cartridge side terminal group between when the cartridge is inserted into the cartridge mounting unit and when the cartridge is mounted to the cartridge mounting unit. There may be cases where the apparatus side terminal group rubs against the side surface of the recessed portion or the like, the constituent member of the cartridge is cut, and shavings are generated. In addition, there may be cases where dust is pulled therein due to static electricity generated by rubbing of the apparatus side terminal group against the side surface of the recessed portion or the like. When shavings, dust, and the like are present at the periphery of the apparatus side terminal group and the cartridge side terminal group, there may be cases where contact between the apparatus side terminal group and the cartridge side terminal group is not properly achieved.\nIn addition, in the techniques of JP-T-2002-505212 and JP-T-2002-513341, the connector is configured to be movable in the X axis direction which is a direction perpendicular to the insertion direction of the cartridge. In addition, in the techniques of JP-T-2002-505212 and JP-T-2002-513341, the recessed portion is provided in the front surface of the cartridge, and the cartridge side terminal group is provided on the bottom surface that is perpendicular to the insertion direction of the recessed portion. In the techniques of JP-T-2002-505212 and JP-T-2002-513341, as a guiding member provided in the recessed portion is inserted into a guiding slot of the cartridge mounting unit, positioning of the cartridge side terminal group and the apparatus side terminal group is performed.\nIn the techniques described above, the cartridge side terminal group is provided on the bottom surface of the recessed portion which is recessed in the reverse direction to the insertion direction and has a relationship of opposing the apparatus side terminal group in the insertion direction. Therefore, when the cartridge is attached to and detached from the cartridge mounting unit, the possibility of the apparatus side terminal group rubbing against the cartridge may be reduced. However, the cartridge side terminal group and the apparatus side terminal group come into point contact with each other between when the cartridge is inserted into the cartridge mounting unit and is connected thereto. Therefore, there is a possibility that contact between the cartridge side terminal group and the apparatus side terminal group may not be achieved due to a slight manufacturing error in the cartridge mounting unit or the cartridge.\nThe problems as described above are significant particularly in a large printer which performs printing on a large-sized sheet (for example, A2 to A0 sizes) because the size of the cartridge is increased.\nThe problems as mentioned above are problems that are not limited to the cartridge and the printer having the cartridge that is detachable, but are common to a printing apparatus that supplies (ejects) different types of printing materials (for example, toner) from ink and a cartridge that is detachably mounted to the printing apparatus."}
{"text": "1. Technical Field\nThe present disclosure relates to portable systems, and particularly to a portable system with a power adapter.\n2. Description of Related Art\nNowadays, most portable devices, such as notebook computers, mobile phones, cameras, and personal digital assistants, display the percentage or level of voltage or charge remaining on the rechargeable batteries. However, for some portable electronic devices, when the rechargeable battery is being charged, and the portable device is not powered on, the level of charge may not be displayed. Therefore, the charging progress of the rechargeable battery cannot be known, which is inconvenient."}
{"text": "1. Field of the Disclosure\nThe present disclosure relates generally to a method for outputting audio signals and an electronic device for supporting the same. In particular, the present disclosure relates to a method and an electronic device for adjusting an audio signal by a filter and outputting the adjusted audio signal.\n2. Description of the Related Art\nAudio data is reproduced by means of a multimedia player to then be output through a speaker. A faithful output of the original sound depends on the performance of the speaker and the characteristics of an audio processing unit of the player. Various techniques have been developed in order to faithfully reproduce the original sound.\nWith recent developments in technology, the audio processing unit may obtain the original sound by using a loudness equalization process that strengthens a low level signal to compensate for the non-linear characteristics of human ears. Another audio processing unit may generate harmonic waves by forming an absolute value by means of the rectifier arrangement. The audio processing unit may process the audio data based on the generated harmonic waves.\nIf audio data is input into two channels, the audio data may be output to speakers that correspond to the two input channels, respectively. Accordingly, there may be limitations on the increase in the low-band performance of the audio data."}
{"text": "Recursive sequence recombination entails performing iterative cycles of recombination and screening or selection to \"tevolve\" individual genes, whole plasmids or viruses, multigene clusters, or even whole genomes (Stemmer, Bio/Technology 13:549-553 (1995)). Such techniques do not require the extensive analysis and computation required by conventional methods for polypeptide engineering. Recursive sequence recombination allows the recombination of large numbers of mutations in a minimum number of selection cycles, in contrast to traditional, pairwise recombination events.\nThus, recursive sequence recombination (RSR) techniques provide particular advantages in that they provide recombination between mutations in any or all of these, thereby providing a very fast way of exploring the manner in which different combinations of mutations can affect a desired result.\nIn some instances, however, structural and/or functional information is available which, although not required for recursive sequence recombination, provides opportunities for modification of the technique. In other instances, selection and/or screening of a large number of recombinants can be costly or time-consuming. A further problem can be the manipulation of large nucleic acid molecules. The instant invention addresses these issues and others."}
{"text": "The present invention relates to a transfer system for carrying out assembling of parts, while circulating a pallet with a work placed thereon along a transfer line forming a closed loop.\nSuch a transfer system for carrying out the assembling of parts, while circulating a pallet with a work placed thereon along a transfer line forming a closed loop is conventionally known (see Japanese Patent Publication No.4-57555). In the above known transfer system, a pallet-shaped carriage on which a work can be placed is circulated on rails placed along a quadrilateral transfer line. In this case, the carriage is allowed to travel by connecting and disconnecting an endless chain, circulatively driven along the rails by a drive source, to or from the carriage, or the carriage is allowed to travel pitch by pitch by bringing the carriage into and out of engagement with a rack which is reciprocally driven through a predetermined distance along the rails by a drive source, and during this time, the assembling of parts to the work placed on the carriage is carried out.\nHowever, the above known system suffers from a problem that when the length of the transfer line is to be changed, it is necessary to change the length of the endless chain or of the rack provided along the transfer line, and such replacing work becomes of a large scale, resulting in a great deal of cost required.\nMoreover, the known system also suffers from another problem that it is necessary to supply parts to be assembled to a work from the outside, to the carriage in succession and for this reason, it is difficult to enhance the working efficiency. To carry out the operation of assembling the parts to the work, working air is required for enabling lifting and lowering of the work to a level which makes the operation easy, or for driving an air tool such as an impact wrench. However, there is a problem that if the working air is supplied from the outside of the moving carriage, the handling or the attachment and detachment of an air hose is troublesome.\nThe present invention has been accomplished in view of the above circumstances, and it is a first object of the present invention to provide a transfer system for carrying out the assembling of parts, while circulating a pallet having a work placed thereon along a transfer line forming a closed loop, wherein the length of the transfer line can be changed to any other one without modification of a pallet driving device.\nIt is a second object of the present invention to provide a transfer system for carrying out the assembling of parts, while circulating a pallet having a work placed thereon along a transfer line forming a closed loop, wherein the operation of assembling parts to the work can be carried out efficiently.\nTo achieve the first object, according to a first aspect and feature of the present invention, there is provided a transfer system having transfer-in stations and transfer-out stations for assembling parts to a work on a pallet, while circulating the pallet along a transfer line forming a closed loop, the transfer system comprising a first conveyer and a second conveyer which are disposed in parallel to each other to transfer a plurality of pallets in an abutting state, a first traverser for transporting the pallets individually from a terminal end of the first conveyer to a start end of the second conveyer, and a second traverser for transporting the pallets individually from a terminal end of the second conveyer to a start end of the first conveyer, the first and second conveyers and the first and second traversers being disposed in the same horizontal plane to form the transfer line, the transfer-in stations and the transfer-out stations being disposed adjacent each other on or in the vicinity of the first and second traversers, the first and second conveyers being operable to drive the pallets by bringing friction rollers which are rotated by a drive source mounted at a fixed side, into abutment against the pallet.\nWith the above arrangement, the first conveyer, the first traverser, the second conveyer and the second traverser are disposed in the same horizontal plane to define the transfer line forming the closed loop for transferring the pallets, and the first and second conveyers drive the pallets by bringing the rotated friction rollers into abutment against the pallet. Therefore, the number of the pallets can be increased and decreased only by changing the lengths of the first and second conveyers without making any modification to the friction rollers, and thus it is extremely easy to change the length of the transfer system. Moreover, while keeping the structure simple and of a low cost, the generation of vibration and noises can be prevented to drive the pallet smoothly. In addition, the transfer-in stations and the transfer-out stations for carrying out the transfer of the work onto and out of the pallet are disposed adjacent each other on or in the vicinity of the first and second traversers. Therefore, most of the transfer line is used ensuring space and time for carrying out the assembling of the parts to the maximum, along which line the work can be transferred, whereby useless space and time required for returning of an empty pallet can be maintained to the minimum.\nTo achieve the first object, according to a second aspect and feature of the present invention, in addition to the first feature, the transfer-in station for transferring the work onto the pallet and the transfer-out station for transferring the work from the pallet are disposed adjacent to each other on or in the vicinity of the first and second traversers.\nWith the above arrangement, the transfer-in and out stations for transferring the work onto and out of the pallet are disposed adjacent each other on or in the vicinity of the first and second traversers. Therefore, most of the transfer line is used ensuring space and time for carrying out the assembling of the parts to the maximum, along which line the work can be transferred, whereby useless space and time required for returning of an empty pallet can be maintained to the minimum.\nTo achieve the first object, according to a third aspect and feature of the present invention, in addition to the first feature, the transfer-in station for transferring a part tray to the pallet and the transfer-out station for transferring the part tray from the pallet are disposed adjacent each other on or in the vicinity of the first and second traversers.\nWith the above arrangement, the transfer-in and out stations for transferring the part tray to and out of the pallet are disposed adjacent each other on or in the vicinity of the first and second traversers. Therefore, most of the transfer line is used ensuring space and time for carrying out the assembling of the parts to the maximum, along which line the part tray can be transferred, whereby useless space and time required for returning of an empty pallet can be maintained to the minimum.\nTo achieve the first object, according to a fourth aspect and feature of the present invention, in addition to the first feature, the first and second traversers include transfer means for transporting the work at a terminal end of either one of the first and second conveyers to a start end of the other of the first and second conveyers.\nWith the above arrangement, the work at the terminal end of either one of the first and second conveyers can be transported to the start end of the other of the first and second conveyers by the transfer means provided on the traversers and hence, the smooth transporting of the pallet can be carried out.\nTo achieve the first object, according to a fifth aspect and feature of the present invention, in addition to the fourth feature, the transfer means have transfer surfaces capable of being lifted and lowered relative to the bottom surface of the pallet on the first and second conveyers, between positions which are higher and lower respectively than the bottom surface of the pallet.\nWith the above arrangement, the transfer surface of the transfer means is lifted and lowered relative to the bottom surface of the pallet on the first and second conveyers, between positions which are higher and lower respectively than the bottom surface of the pallet. Therefore, the transporting of the pallet between the first and second traversers and the first and second conveyers can be conducted smoothly.\nTo achieve the first object, according to a sixth aspect and feature of the present invention, in addition to the fourth feature, the transfer system further includes a cylinder for pushing out the pallet from the terminal ends of the first and second traversers to the start ends of the first and second conveyers.\nWith the above arrangement, the pallet is pushed out from the terminal ends of the first and second traversers to the start ends of the first and second conveyers by the cylinder and therefore, the transporting of the pallet from the first and second traversers to the first and second conveyers can be conducted smoothly.\nTo achieve the first object, according to a seventh aspect and feature of the present invention, in addition to the first feature, the transfer system further includes an air tank which is mounted on the pallet and in which air for driving an air tool is stored, and an air supply means provided at least at one place on the transfer line for supplying air to the air tank.\nWith the above arrangement, the air tool can be driven by the air stored in the air tank mounted on the pallet and therefore, the handling of an air hose is simplified, and also the operability of the air tool is improved, as compared with a case where air is supplied from the outside of the pallet. Moreover, the air can be supplied to the air tank, while the pallet is being circulated along the transfer line. Therefore, the supplying of the air to the air tank can be conducted easily.\nTo achieve the second object, according to an eighth aspect and feature of the present invention, there is provided a transfer system having a transfer-in station and a transfer-out station for assembling parts to a work on a pallet, while circulating the pallet along a transfer line forming a closed loop, the pallet including a part support for supporting parts and an air tank in which working air for assembling the parts to the work is stored, and an air supply means being provided at least at one place on the transfer line for charging air in the air tank.\nWith the above arrangement, the pallet includes the part support for supporting the parts to be assembled to the work and hence, it is unnecessary to sequentially supply the parts from the outside to the pallet moving along the transfer line, leading to an enhanced operability. Moreover, the parts can be assembled to the work by the air stored in the air tank mounted on the pallet without sequential supplying of the air from the outside to the pallet moving along the transfer line, and hence, the handling of an air hose is simplified. In addition, the air can be supplied to the air tank, while the pallet is being circulated along the transfer line. Therefore, the supplying of the air to the air tank can be conducted easily.\nTo achieve the second object, according to a ninth aspect and feature of the present invention, in addition to the eighth feature, the work supported on work support jigs provided on the pallet is lifted and lowered by air cylinders connected to the air tank.\nWith the arrangement, the work supported on the work support jigs is lifted and lowered by the air cylinders and hence, the work can be supported in a position suitable for the operation, whereby the part assembling operation can be carried out efficiently. Moreover, the air cylinders can be driven by the air stored in the air tank mounted on the pallet without sequential supplying of the air from the outside to the pallet moved along the transfer line and hence, the handling of an air hose to the air cylinders is simplified.\nWith the above arrangement, the air tool is driven by the air stored in the air tank mounted on the pallet without sequential supplying of the air from the outside to the pallet moved along the transfer line. Therefore, the handling of the air hose to the air tool is simplified, and also the operability of the air tool is improved.\nTo achieve the second object, according to an eleventh aspect and feature of the present invention, in addition to the eighth feature, the work support jigs provided on the pallet are capable of being lifted and lowered to levels suitable for the operation.\nWith the above arrangement, since the work support jigs provided on the pallet are capable of being lifted and lowered to the level suitable for the operation, the work can be supported in a position suitable for the operation, whereby the part assembling operation can be carried out efficiently.\nTo achieve the first object, according to a twelfth aspect and feature of the present invention, in addition to the first feature, a working space, in which an operator can move, is provided inside the transfer line forming the closed loop, so that access of the operator to the pallet is enhanced by the working space.\nWith the above arrangement, it is possible for the operator to access any of the pallets through the working space at a low speed. The brake rollers 23, 23 apply a braking force to one of the plurality of pallets P moved by the main drive rollers 22, 22, which is located at the terminal end of the first conveyer C1, thereby preventing the pallet P from overrunning onto the first traverser T1. The transfer-out drive rollers 24, 24 are driven at a speed higher than that of the main drive rollers 22, 22, when one pallet P located at the terminal end of the first conveyer C1 is to be separated from the succeeding pallet P and moved to the start end of the first traverser T1. The transfer-in drive rollers 21, 21 are driven at a speed higher than that of the main drive rollers 22, 22, when one pallet P located at the terminal end of the second traverser T2 is to be moved to the start end of the first conveyer C1 to catch up with the final pallet P on the first conveyer C1."}
{"text": "In the magnetic tapes for the audio, video or computer purpose, in general, there is recently an increasing demand for high-density recording, so that the particle size of the conventional Co-containing magnetic iron oxides becomes finer.\nAttention is also paid to the ferromagnetic metal thin-film type magnetic recording media, which rely upon ferromagnetic alloy powders, or are formed by the vapor deposition techniques such as vacuum deposition, sputtering and ion plating, and such media have been put to practical use.\nIn such audio, video or computer magnetic tapes, the magnetic layers are finished to a certain smoothness to improve the sensitivity (especially output in a high-frequency range). In some cases, however, improvements in sensitivity could not be introduced even with the magnetic layers having a smoothly finished surface, since the irregularity of the backing layers is transferred thereto.\nIn recent years, the total thickness of magnetic tapes tends to be reduced as much as possible, and is in most cases no higher than about 20 microns.\nThis is accompanied by a reduction in the thickness of the associated backing layers, and often causes that they are coated in a thickness of no higher than 2 microns. For that reason, wear of the backing layers is likely to occur.\nThus, the reduction in the total thickness of the magnet tapes is accompanied by a drop in the strength thereof, and have adverse influences upon running characteristics, running durability, taking-up characteristic, output variations, etc."}
{"text": "Among the most useful developments in modern times is the use of electronic mail or messages (email) for efficient correspondence around the world. As numerous email items accumulate in the inbox or data store for a user, attempts are made to sort the items by priority, follow-up, alphabetically, or by filing the emails in folders for future reference or follow-up. Manually filing emails in folders can be an undesirable and costly activity yielding fruitless results. Frequently emails sought for retrieval cannot be found easily or at all because a user may not remember where the email is filed or located.\nPrevious email applications have the capability to auto-file and search for email items meeting user specified search criteria, however auto-file rules scatter a user's messages to different folders before they have been viewed. Furthermore, the search filters have complicated dialogs and the searches are time consuming and costly intensive and time consuming thereby hindering the utility of existing email applications. This problem is further complicated by the incapability of email items to exist in more than one folder without copying the item to multiple folders. When email items are copied to multiple folders and one copy is updated or deleted the other copies are stored in their original form thereby leading to confusion and a waste of memory.\nAccordingly, there is a need for a method, system, and computer-readable medium for providing persisting and continuously updating search folders that allow for efficient, consistent, and user-friendly organization and reorganization of email items. It is with respect to these and other considerations that the present invention has been made."}
{"text": "This invention relates to a leak detection system and process for detecting and ejecting defective containers. Particularly, this invention relates to a leak detection system for filled and sealed pliable containers.\nThe leak detection system of this invention is used to test filled and sealed containers or packages for leakage and to automatically eject any defective containers from the product stream. Leaks in plastic molded containers, such as plastic bottles, can result from impurities in the plastic materials or in the molding process of the containers. Typically, leaks can be large, to cause fluid or other product contents to pour from the containers, or the leaks can be small, to cause small quantities of air or product contents to be expelled from the container particularly when placed under pressure.\nDue to advances in the manufacture and filling of pliable or plastic containers, wherein the containers are filled and sealed in a continuous process, as well as in the speed of these processes, it has become increasingly necessary to properly test the individual containers for product leakage. Subsequent to the manufacturing, filling and sealing processes, the containers may be stacked or otherwise grouped, and/or packaged, for example, in cardboard boxes. Leaking containers can cause sanitation problems and may cause contamination in the packaged containers. Any such leakage is not desirable, particularly in the food services and pharmaceutical industries. Moreover, leaking containers cause warehousing or storage problems in situations where containers, particularly those having liquid contents, are stacked in vertical layers. Thus, it is desirable to properly test any filled container prior to further packaging.\nMost prior art devices and leak detection processes are designed for testing individual product containers or bottles in bath-type operations and, therefore, are not well suited for use in continuous product processing operations. Thus, because of the continuous nature and advancement in these manufacturing processes, most prior art leak detection systems proposed and utilized in the past are not suited for such operations.\nThe leak detection system of this invention, overcomes the shortcomings and limitations of these prior art devices. particularly, the leak detection system of this invention permits filled and sealed pliable containers to be tested in a continuous operation."}
{"text": "Thrust vectoring is an essential feature for vehicles capable of transitioning in flight direction and/or transitioning between various types of flight modes, such as from a takeoff mode to a forward flight mode. In particular, propulsors (e.g., propellers, rotors, fans, ducted fans, or other thrust generating devices that accelerate air) may be designed to rotate or tilt about the vehicle to accomplish a transition in flight direction. Thrust vectoring devices provide a mechanism by which such propulsors move relative to the vehicle. Thrust vectoring devices may be employed in a number of vehicle types, including but not limited to vehicles configured for vertical takeoff and landing (“VTOL”), conventional takeoff and landing (“CTOL”), short takeoff and landing (“STOL”), and the like. Current thrust vectoring devices employ mechanically driven actuators to facilitate movement of the propulsors. With these current systems, a substantial amount of machinery, hardware and/or software may be required to appropriately engage the thrust vectoring, and often must work against freestream velocities around the vehicle to accomplish a change in flight mode and/or flight direction.\nPrior art solutions for thrust vectoring system have not resolved the need for an approach to perform one or more of the above actions without drawbacks, e.g., mechanical/electrical complexity, size and weight constraints, and/or cost-prohibitive. Therefore, there is a need for thrust vectoring systems and methods that address one or more of the deficiencies described above amongst others."}
{"text": "1. Field of the Invention\nThe present invention relates to methods and systems for constructing XML query, more particularly, to methods and systems for constructing XML query to schema variable XML documents.\n2. Description of the Related Art\nXML (eXtensible Markup Language) has been widely adopted in many domains to store and exchange data. In some domains, for example, the CDA (Clinical Document Architecture) for the healthcare industry or the XBRL (extensible Business Reporting Language) for the financial industry, the variable schema, which is a set of XML documents that share a similar but not an identical XML schema, is a common feature that makes it difficult for users to consume data in XML documents.\nIn order to construct a correct XQuery/XPath to extract data from these XML documents, users must completely understand each schema and how each data element in each schema relates to other elements, that is, the relationship of XML elements in a same XML document tree and the relationship of different XML trees, which usually takes too much effort or in some extreme circumstances (e.g., when there are too many schemas) may be impossible for users.\nMDA (Model-Driven-Approach) modeling methods are used in many fields in order to use unified syntax to express and exchange complex data in a flexible manner. Using this methodology, by defining and modeling semantics, vocabularies, data structures and data at levels of meta model, model (schema) and instance respectively, different users can make extension/derivation from basic schemas according to their different requirements, so that various data express and exchange requirements can be satisfied under different circumstances.\nFIG. 1 shows an XML modeling architecture which follows the MDA methodology 1000. The top level of the architecture is a Meta Model layer 1010, which includes the definitions of syntax and semantics of modeling notations, a Model level 1020 is defined according to the Meta model, which includes a terminology model 1021 defining common vocabularies and terms of specific domain, and a schema model 1022 defining common data structures of how to organize the common vocabularies and terms. At the lowest level, Data 1030, i.e., Instance layer, is provided under the Model level 1020.\nFIG. 2 describes the hierarchical MDA modeling methodology shown in FIG. 1 using the XBRL as an example. As shown in FIG. 2, in the XBRL specification defined by the XBRL international organization, a Meta model is defined, including notations that can be used by financial reports to tag reporting items. For example, the ‘substitutionGroup’ attribute of each reporting item in a financial report should be tagged as either “item” or “tuple”. Based on the meta model, regulators such as the SEC (Securities Exchange Commission) have defined some basic schemas and terms, e.g., vocabularies of reporting items. For instance, the SEC has defined the reporting item “Revenue” to be used in financial reports, and “Revenue” elements can be represented in income statements. Companies can extend these basic schemas by themselves to customize their own reporting templates. For example, if Company A has a definition in its schema extension, “Income”-“Revenue”, then in the income statement of Company A, element “revenue” should be represented as sub node of element “income”. If Company B has a definition in its schema extension “Accrued Income”-“Revenue”, then this indicates that element “Revenue” is sub node of element “Accrued Income” in the income statement of Company B. Each company can also have specific data facts according to the extended schemas, e.g., specific data of reporting items defined in these schemas. For example, the value of the reporting item “Revenue” in the financial report of company A is “1,000,000 USD”, the value of the item “Revenue” in the report of company B is “1,500,000 USD”.\nFIG. 3 describes the hierarchical MDA modeling architecture shown in FIG. 1 using the CDA as another example. As shown in FIG. 3, the CDA international organization has defined the syntax and semantics of notations that can be used by electronic clinical documents, e.g., a clinical document include data for specific elements such as “Entity”, “Observation”, “Symptom”, “Body Structure”, etc. Based on the above syntax, a standard organization, SNOMED (Systematized Nomenclature of Medicine) has further defined some common terms, i.e., vocabularies of clinical data. For example, SNOMED has defined terms, such as “Shadow” for “Observation” and “Chest” for “Body Structure”. Based on such common syntax and terms, hospitals and equipment manufacturers can create extensions on the basic schema to obtain their own CDA document schemas. For example, hospital A can define a nested structure like “Entity”-“observation”-“body”-“symptom” and hospital B can define a nested structure like “Entity”-“body structure”-“symptom”. Hospital A can produce data facts like “Tom”-“SNOMED CT”-“Chest”-“Shadow”, indicating “a shadow of Tom's chest observed with SNOMED CT”, which corresponds to the above schema: “Entity”-“observation”-“body”-“symptom”; and hospital B may produce data facts like “Lee”-“Chest”-“Shadow”, indicating a shadow of Lee's chest, which corresponds to the above schema “Entity”-“body structure”-“symptom”.\nUsers can query data in above XML documents. As shown in FIG. 2, financial data consumer (e.g., bank or investor) may want to know “the revenue changing trends of companies A and B in the past three years.” Similarly, as shown in FIG. 3, a clinic data consumer (e.g., doctors or medicine makers) may want to know “the list of patients with shadow in chest.” However, under existing techniques in the prior art, it is difficult to formulate queries because even though XML documents share a common syntax and vocabulary model, they can still be based on different schemas or data structures. When constructing an XML query for a document, even though there is a common vocabulary model that covers all elements involved in the queries, users still need to know the specific schema of each instance document. Data consumers have to understand each extended schema in order to correctly the XQuery/XPath thereby imposing a heavy burden on users.\nAn XQuery/XPath is a known query language in the prior art. XPath is a language for selecting nodes from XML documents. XQuery is a query language (with some programming language features) designed to query on XML data sets.\nReferring to FIG. 2, even though the financial report of companies A and B share the common vocabularies for income statements (e.g., “Revenue”), they have different data structure to organize these same elements because of their own schema extensions. When querying the “Revenue” of companies A and B, the specific schema structures have to be understood by a user, i.e., in the extended schema defined by company A, element “Revenue” is the child node of element “Income”, and in the extended schema defined by company B, element “Revenue” the child node of element “Accrued Income”. It is necessary for users to get above information to construct a correct XPath/XQuery.\nReferring to FIG. 3, although sharing the common syntax defined by CDA and the vocabularies defined by SNOMED, it is still possible for hospitals to extend the basic schema to obtain customized schemas that are suitable to their own business requirements. Examples include different customized schema like “Entity”-“observation”-“body”-“symptom” and “Entity”-“body structure”-“symptom”. When answering questions like “patients with shadow in chest”, information about detail structures of each different schema is required. For example, with respect to patient “Tom”, the element “Shadow” is sub node of element “Chest”, which is sub node of element “Observation”, which in turn is the sub node of root element “Tom”. As to patient “Lee”, the element “Shadow” is sub node of element “Chest”, which in turn is the sub node of root element “Lee”. Accordingly, if there are a number of extended schemas, users need to know the specific structures of each schema in order to construct queries even if those schemas share the common syntax and vocabularies, i.e., the basic schema. This is an unnecessary burden to users.\nFIG. 4 shows a set of shared vocabularies, including “element 1” to “element 5”. These elements are referred by the extended schemas shown in FIGS. 5a and 5b, but they are organized in element trees with different structure in the extended schemas respectively. Comparing the different schemas shown in FIG. 5a and FIG. 5b, it can be seen that information of specific element nesting paths is required even though the same elements are queried. Therefore, in the prior art, when querying on XML documents with these extended schemas, for example, when querying elements 3, 4 and 5, paths to the elements being queried in the corresponding XML documents are first determined before constructing queries according to those paths. As mentioned before, this approach has a significant drawback because users have to know the specific structures of each schema.\nAnother existing approach in the prior art is a schema-less query by wildcard. However, wildcard queries can lose the relationship between elements being queried, and therefore is limited to queries for a single element. In the case where multiple elements are involved and the relationship therebetween has to be considered, wildcard queries cannot return desired results.\nFIG. 6 shows the above shortcoming of wildcard query. A simple XML schema is shown on the top of FIG. 6, where “element—1” (sub node) and “element—2” (sub node) are located under “element_a” (root node).\nAn XML document is provided in FIG. 6:\n<element_a>  <element_1>1-1</element_1>  <element_2>1-2</element_2></element_a><element_a>  <element_1>2-1</element_1>  <element_2>2-2</element_2></element_a>\nThe above XML segment has two data facts for “element_a”, each of which includes two data facts “1-1”, “1-2” and “2-1”, “2-2” for “element—1” and “element—2” respectively. Assume a user want to get all of data facts of “element—1” and “element—2”. When making a wildcard query (that is, a schema-less query), the user inputs query “\\\\element—1” and “\\\\element—2”, and gets query results shown at the bottom of FIG. 6. Although data facts “1-1”, “1-2”, “2-1”, and “2-2” are returned, the fact that the data facts “1-1”, “1-2” should be joint and data facts “2-1”, “2-2” should be joint are lost. In fact, what user really expects is the table shown in the middle of FIG. 6, wherein data facts “1-1”, “1-2” and data facts “2-1”, “2-2” are associated with each other respectively.\nTherefore, existing approaches do not support how to make a query of data in schema variable XML documents even users do not have any knowledge about the data structures of specific XML schema, while the relationship between the elements are maintained in the returned query results."}
{"text": "This disclosure relates generally to finishing a component surface and, more particularly, to finishing techniques that immerse at least a portion of the component in a fluid.\nMany components include surfaces that are finished prior to use. The surfaces of both used components and newly manufacturing components may be finished. As an example, the surfaces of a high pressure compressor blade may be polished prior to use within a turbomachine. Polishing is one type of finishing. Turbomachine components often require relatively low surface finishing, such as an average surface roughness (Ra) of 1 to 5 micro-inch (0.0254 to 0.127 microns) Ra.\nElectrochemical machining (or ECM) is an example finishing process that immerses at least a portion of a component in a fluid. Electrochemical machining processes use electrical energy to remove material. An ECM process creates an electrolytic cell in an electrolyte medium. A tool serves as a cathode, and the component serves as an anode. During the process, the ECM tool is positioned very close to the workpiece and a low voltage, high amperage DC current is passed between the two via an electrolyte. A high-amperage, low-voltage current is then applied to dissolve and remove metal from the component.\nElectrochemical deburring is a type of electrochemical machining designed to remove burrs and impart small radii to corners. The process normally uses a specially shaped electrode to carefully control the process to a specific area. Unlike many other machining methods, workpiece hardness is not a significant factor in ECM, making ECM suitable for difficult-to-machine materials.\nThe electrochemical machining process can produce byproducts, such as salt, at the surface of the component being finished. The salt inhibits the finishing."}
{"text": "As is known in the art, conventional SATCOM terminals utilize small or low profile reflector antennas in applications having significant size constraints. The small or low profile reflector antennas typically include a feed assembly that transmits signals from a transmitter or receives signals to a receiver in the respective antenna system. However, the size of the feed assembly can limit the type of SATCOM applications the small or low profile reflector antennas can be utilized in. Further, many feed assemblies are only configured to provide and support single-band or dual-band operation."}
{"text": "A hot-dip galvanized steel sheet has been frequently used for automobiles, home electric appliances, building materials, and the like. Extremely high formability is required for an automotive steel sheet which is pressed into complicated shapes. In addition, since a demand for rust-proof performance of the automobile increases in recent years, the hot-dip galvanized steel sheet comes to be often used for the automotive steel sheet.\nIn recent years, there are increasing needs for a high-strength steel sheet excellent in strength and ductility in view of reduction in vehicle body weight. For example, Patent Document 1 discloses a steel sheet having a steel sheet structure where a ferrite phase, a bainite phase, and an austenite phase are mixed. In addition, it is disclosed that this steel sheet is a steel sheet making use of transformation-induced plasticity exerting high ductility by the transformation of retained austenite into martensite at a forming time.\nThis type of steel sheet contains, for example, 0.05 to 0.4 mass % of C, 0.2 to 3.0 mass % of Si, and 0.1 to 2.5 mass % of Mn, and forms a complex structure by controlling a temperature pattern during a cooling process after annealing at a two-phase region. There is a characteristic that desired properties can be secured without using expensive alloying elements. In recent years, there are increasing needs for a high-strength hot-dip galvanized steel sheet where a steel sheet base material surface is subjected to hot-dip galvanized plating in order to secure the rust-proof performance even for the aforementioned high-strength steel sheet.\nThere are increasing cases where the high-strength steel sheet is used not only as a strengthening member for an inner plate use but also as an outer surface member which may receive impact from flying stones or obstacles at a vehicle body driving time. When the high-strength steel sheet is applied to a complicated shaped member, high processability is required. The high-strength hot-dip galvanized steel sheet is required to secure high plating adhesiveness evaluated by severe evaluation methods such as a ball impact test and a draw bead test in addition to a normal 60° V bending test assuming cases when impacts from flying stones or obstacles are received during driving and the plating adhesiveness at a hard working time.\nSince cracks occur at plating and a base iron regarding a portion receiving extremely severe working such as a 180° bending work vertex part, corrosion is likely to occur from the portion even after it is subjected to a conversion treatment and an electrodeposition coating. Even a slight corrosion may cause hydrogen penetration from the corroded part to thereby increase a possibility of hydrogen embrittlement cracking particularly when the base material is the high-strength steel sheet.\nWhen the high-strength steel sheet is subjected to galvanized plating at a continuous hot-dip galvanizing facility, plating wettability largely decreases when an Si amount of the steel sheet exceeds 0.3 mass %. There is therefore a problem that unplating occurs to deteriorate an appearance quality in case of a Sendzimir method using a normal Al containing plating bath. In addition, the plating adhesiveness at the impact time or hard working time was difficult to be secured because the plating adhesiveness is simultaneously largely lowered.\nIt is said that this is because an external oxide film which contains oxides containing Si and Mn having low wettability with respect to molten Zn is generated at a steel sheet surface at a reduction annealing time.\nAs a means to solve this problem, Patent Document 2 proposes a method where the steel sheet is previously heated in an atmosphere at an air ratio of 0.9 to 1.2 to generate Fe oxide, and then plating is performed in a bath where Mn and Al are added after a thickness of the oxide is set to 500 Å or less in a reduction zone containing H2. However, there have been problems that it is difficult to accurately control the thickness of the oxide, and a manufacturing conditional range in an actual machine is narrow because various steel sheets containing various additive elements are passed through an actual line. Further, though an effect of improving wettability and plating adhesiveness at a normal working time could be expected, an effect of improving the plating adhesiveness at the impact time or hard working time was small.\nPatent Document 3 discloses a method improving the platability by supplying specific plating to a lower layer as another unplating suppression means. However, this method requires to newly provide a plating facility at a previous stage of an annealing furnace in a hot-dip plating line, or previously perform a plating treatment in an electroplating line. Both cases cause serious increase in manufacturing cost. Further, the effects of improving the plating adhesiveness at the impact time or hard working time and the worked portion corrosion resistance were small.\nMeanwhile, Patent Document 4 discloses a method manufacturing a hot-dip galvanized steel sheet without oxidizing Fe in the steel sheet by adjusting an oxygen potential in an annealing atmosphere at the annealing time. In this method, easily oxidizable elements such as Si and Mn in the steel are internally oxidized by controlling the oxygen potential of the atmosphere to suppress formation of an external oxide film, to thereby improve the platability. Although sufficient adhesiveness can be secured at the normal working time by applying this method, the effects of improving the plating adhesiveness at the impact time or hard working time and the worked portion corrosion resistance could not be expected."}
{"text": "Integrated circuits are available in many different packages, technologies, and sizes. Most integrated circuits are available in plastic packages, which are generally intended for commercial operating environments at a low cost. Commercial operating environments have a specified operating range from 0° C. to 70° C. Integrated circuits for military applications have historically been packaged in either metal or ceramic hermetic packages, which are able to work reliably in more demanding environments than commercial integrated circuits. Military operating environments have a specified operating range from −55° C. to 125° C. In order to save costs, the military has purchased integrated circuits through COTS (Commercial Off-The-Shelf) programs. However, these components are generally commercial grade components in plastic packages, and not intended for demanding environments requiring the broader temperature range reliability and durability of ceramic and metal hermetically packaged integrated circuits.\nDepending on size and complexity, integrated circuits are available in a wide range of packages. Although many older integrated circuits were packaged using through-hole technology packages, surface mount packages have dominated over the past several decades. Surface mount packages generally have circuit density, cost, and other advantages over through-hole integrated circuits. Examples of through-hole packages include DIP (dual-in-line plastic) and PGA (pin grid array). Examples of surface mount packages include SOIC (small-outline integrated circuit) and PLCC (plastic leaded chip carrier).\nIntegrated circuit packages generally consist of a semiconductor die placed within a package base and bonded to the base with a suitable die attach adhesive. In conventional technology, the die is electrically attached to a lead frame of the package base with discrete bond wires, which connect individual pads of the die with package leads. In most cases, the bond wires are gold, but in other environments can be copper or aluminum. Specialized equipment is required to attach the bond wires to the die pads the lead frame. Once all of the bond wires are attached, the package lid is bonded to the package base and the integrated circuit can be tested."}
{"text": "2. Origin of the Invention\nThe invention described herein was made in the performance of work under a NASA contract, and is subject to the provisions of Public Law 96-517 (35 USC 202) in which the contractor has elected to retain title.\n2. Field of the Invention\nThe present invention relates, in general to real-time data compression of raw imaging data for high resolution imaging radar systems, and in particular to a block adaptive quantizer and an associated VLSI processor to provide real-time data compression for high resolution imaging radar systems."}
{"text": "Automated or semi-automated directional positioning control of sensors, including, for example, image sensor such as video and still cameras operating over a range of electromagnetic frequencies, is important in a variety of applications. Such control may be utilized in manufacturing tooling and other applications but is more widely known as desirable for various surveillance and security applications, and is specifically important for use with aircraft to capture visual data from above. In one example, surveillance equipment may be used in an Unmanned Aerial Vehicle (UAV) with sensor directional control to provide operator remote sensing and aerial situational assessment capabilities. In other examples, sensors may provide national airspace traffic sense-and-avoid threat detection. In still other examples, sensors may enable guidance of weaponry such as a cruise missile, where either an operator or an autonomous collision avoidance system rely on image sensor feedback to control the missile. In some examples, surveillance or other systems may employ systems or devices that provide for automated or semi-automated control of sensors, for example directional aiming of an image sensor. Such devices may also provide for motion stabilization of sensors.\nSystems and devices that provide for the control and/or stabilization of image sensors may be referred to as gimbals. Gimbals enable the automatic or semi-automatic control of directional sensor positioning about a range of motion via various mechanical, electrical, and/or magnetic mechanisms. The range of motion of a gimbal may be about a single axis, or about multiple axes.\nIt may be advantageous to display and/or process data collected by a visual sensor of a gimbal system in real time, or close to real time. However, such data is often relatively large in size, which may render data transfer from the gimbal difficult.\nA wired data connection to transmit data from a sensor of a gimbal system may cause forces that impede an ability of a gimbal to effectively position a visual sensor. Also, due to frequent and forceful gimbal movement, a wired data connection may become physically damaged over time, rendering a gimbal system wholly or partially inoperative.\nSome gimbal systems employ one or more mechanical slip rings to transfer data from a sensor of the gimbal system. The mechanical slip rings may be placed on each axis of motion of a particular gimbal system. Aspects of a mechanical slip ring may limit an ability of designers to minimize the size of the gimbal system. Electromagnetic characteristics of mechanical slip rings may further cause the slip rings to have a limited data transfer rate."}
{"text": "Wire bonding is typically applied to make electrical connections between an integrated circuit die or chip and a carrier on which the die is mounted. Bonding wires are attached to bond pads on the chip and bonding leads on the carrier respectively by ultrasonic welding using an ultrasonic transducer which is integrated into a wire bonding apparatus. The ultrasonic transducer is an energy converting-device which converts electrical energy into ultrasonic vibrations and transmits the ultrasonic vibrations to a capillary at a tip end of the transducer to perform wire bonding.\nFIG. 1 illustrates a side view of a conventional wire bonding apparatus 10 with a transducer 12 positioned over a bonding area for wire bonding. The transducer 12 has a long slender rod-shaped horn 13 and a capillary 16 located at a tip end of the horn 13. Wire is threaded through the capillary 16 and extends out of the capillary tip. The capillary 16 is raised or lowered with respect to a bonding area by moving the transducer 12 vertically.\nAn ultrasonic generator 17 is connected to an end of the horn 13 that is opposite to the tip end of the horn 13 where the capillary 16 is mounted. The ultrasonic generator 17, typically comprising a stack of piezoelectric elements, is housed in a transducer holder 14 mounted on a bond arm 18. The transducer holder 14 is supported by the bond arm 18, which is in turn connected to a sliding bar 22 at a pivot 20. Rotational motion of the bond arm 18 about the pivot 20 rotates the bond arm 18 up and down thereby raising or lowering the transducer 12 relative to the bonding area. The sliding bar 22 may further move and position the bond arm 18 and transducer 12 along the x and y axes.\nA carrier 26, on which components such as semiconductor dice and bonding wires are attached, is held in position by a window clamp 24 on a wire bonding platform to perform wire bonding. The wire bonding platform includes a heater block 28 which provides heat to the carrier 26 to facilitate wire bonding conducted on it.\nDuring wire bonding, the ultrasonic generator 17 produces ultrasonic vibrations which are transmitted along the horn 13. The long slender rod shape of the horn 13 is suitable for amplifying the ultrasonic vibrations transmitted to the capillary 16 while suppressing attenuation. In this way, ultrasonic vibrations may be transferred to the capillary 16 efficiently.\nA pressing force is also applied to the capillary 16, and bonding is accomplished by applying the ultrasonic energy transmitted through the horn 13 onto the wire that is subsequently attached to bonding pads, typically on the semiconductor dice. Control of the positioning of the capillary 16 relative to the bonding area is essential to perform wire bonding accurately. This is particularly so when very fine wires and bond pitches are involved.\nAs can be seen in FIG. 1, the transducer 12 is positioned over the heating zone during wire bonding. Heat is transmitted to the immediately vicinity of the heater block 28, which means that the transducer 12 is exposed to heat from the heater block 28 by radiation. The heat may cause the transducer 12 and the transducer holder 14 to expand. This may therefore affect the relative positioning of the capillary 16 and a bonding position. As a result, wire bonding accuracy is affected.\nA thermal shield may be used for reducing heat transmission from the heater block 28 to the transducer 12 and the transducer holder 14. FIG. 2 is a side view of the wire bonding apparatus 10 with a thermal shield 30 affixed to the sliding bar 22 of the wire bonding apparatus. The thermal shield 30 is operative to separate the transducer 12 from the heater block 28, thereby reducing radiation of heat from the heater block 28 directly to the transducer 12. However, the window clamp 24 needs to be lifted to unclamp the carrier 26 every time the carrier 26 is moved or when it is removed from the heater block 28. The presence of the thermal shield 30 limits the vertical distance A that is movable by the window clamp 24. It is desirable to increase the vertical distance that is movable by the window clamp 24 to ensure sufficient clearance of the window clamp from the carrier 26 and the components attached thereon when the carrier 26 needs to be moved.\nFIG. 3 illustrates a front view of the wire bonding apparatus of FIG. 2 showing the thermal shield 30 between the transducer 12 and the heater block 28. While the thermal shield 30 may reduce the heat radiated from the heater block 28 to the transducer 12, it is noted that the transducer 12 is still exposed to ambient heat from the surrounding air above the thermal shield 30. The transducer 12 and transducer holder 14 would still expand such as to change the relative positioning of the capillary 16 and the carrier 26, resulting in decreased wire bonding accuracy.\nFIG. 4 is a side view of the wire bonding apparatus 10 of FIG. 2 illustrating movement of the transducer 12 along the y and z axes. The thermal shield 30 is movable together with the sliding bar 22 along the x and y axes as it is attached to the sliding bar 22. However, when the transducer 12 moves along the z-axis by rotary movement of the bond arm 18 about the pivot 20, the thermal shield 30 cannot move correspondingly in the z axis. Thus, the transducer 12 is likely to receive more heat energy from the surrounding air as it moves further away from the thermal shield 30. For the aforesaid reasons, use of a conventional thermal shield in a wire bonding system does not effectively address the problem of thermal instability encountered by a transducer during wire bonding. It would be desirable to provide better thermal protection to the transducer 12 from the heater block 28."}
{"text": "The invention relates to reducing NOx emissions from furnaces by addition of a water biofuel or waste fuel slurry to the furnace outside the primary combustion zone.\nDuring combustion of fuels with fixed nitrogen such as coal, oxygen from the air may combine with the nitrogen to produce nitrogen oxides (NOx). At sufficiently high temperatures, oxygen reacts directly with atmospheric nitrogen to form NOx. Emission of nitrogen oxide is regarded as undesirable because the presence of nitrogen oxide in a furnace flue gas (along with sulfur dioxides) causes the condensed gases to become corrosive and acidic. There are numerous government regulations which limit the amount of nitrogen oxide which may be emitted from a combustion furnace. Titles I and IV of the Clean Air Act as amended in 1990 (xe2x80x9cThe Clean Air Actxe2x80x9d) require significant NOx reduction from large power plants. Title I of the Clean Air Act focuses on the problems of ozone non-attainment. Ozone is formed as a result of photochemical reactions between nitrogen oxides emitted from central power generating stations, vehicles and other stationary sources, and volatile organic compounds. Ozone is harmful to human health. Consequently, in many urban areas the Title I NOx. controls are more stringent than the Title IV limits. Thus, there is a need for apparatus and processes which reduce the nitrogen oxide emissions in furnace flue gas.\nCommercially available techniques to reduce the nitrogen oxide emissions in a furnace flue gas are low NOx burners, overfire air, selective non-catalytic NOx reduction (SNCR), selective catalytic reduction (SCR), and reburning. Currently, retrofitting boilers with low NOx burners and overfire air is the most economic route to comply with Title IV requirements of the Clean Air Act. However, low NOx burners cannot reduce NOx emissions to levels required by Title I of the Clean Air Act. As a consequence, electric utilities are faced with the option of adding SNCR or reburning to the boiler. In addition, cyclone boilers cannot be retrofitted with low NOx burners. SCR, SNCR and reburning are the options for cyclone boilers.\nThe reburning process is also known as in-furnace nitrogen oxide reduction or fuel staging. The standard reburning process has been described in several patents and publications. See for example, xe2x80x9cEnhancing the Use of Coals by Gas Reburning -Sorbent Injection,xe2x80x9d submitted at the Energy and Environmental Research Corporation (EERC), First Industry Panel Meeting, Pittsburgh, Pa., Mar. 15, 1988; xe2x80x9cGR-SI Process Design Studies for Hennepin Unit #1xe2x80x94Project Review,xe2x80x9d Energy and Environmental research Corporation (EERC), submitted at the Project Review Meeting on Jun. 15-16, 1988; xe2x80x9cReduction of Sulfur Trioxide and Nitrogen Oxides by Secondary Fuel Injection,xe2x80x9d Wendt, et al.; published at the Symposium of the Combustion Institute, 1972; xe2x80x9cMitsubishi xe2x80x98MACTxe2x80x99 In-Furnace NOx Removal Process for Steam Generator,xe2x80x9d Sakai, et al.; published at the U.S. xe2x80x94Japan NOx Information Exchange, Tokyo, Japan, May 25-30, 1981. In reburning a fraction of the total thermal input is injected above the primary flame zone in the form of a hydrocarbon fuel such as coal, oil, or gas. A reburn zone stoichiometry of 0.90 (10% excess fuel) is considered optimum for NOx control. Thus, the amount of reburn fuel required is a direct function of the primary zone excess air. Under typical boiler conditions a reburn fuel input in the range of 15% to 25% is sufficient to form a fuel-rich reburn zone. The reburn fuel is injected at high temperatures in order to promote reactions under the overall fuel rich stoichiometry. Typical flue gas temperatures at the injection location are above 2600xc2x0 F. Completion air is added above the fuel rich reburn zone in order to burn off the unburnt hydrocarbons and carbon monoxide (CO). In addition to the above specifications, the prior art on standard reburn teaches that rapid and complete dispersion of the reburn fuel in flue gas is beneficial. Thus, flue gas recirculation (FGR) has been used to promote mixing in all standard reburn demonstrations. Standard reburn technology requires a tall furnace to set up a fuel rich zone followed by a lean zone. Many furnaces do not have the volumes required for retrofitting this technology.\nIn current practice of the reburning process, usually more than enough fuel is added to react with all of the oxygen remaining in the original combustion products. A reducing zone, or a zone with an excess of fuel is formed. In this reducing zone the NO reacts with the excess fuel to form N2, NH3, HCN, and other reduced nitrogen. Then more air is added to combust the remainder of the reburn fuel. At this point the NH3, HCN, and other reduced forms are oxidized to N2 and NO. At this step and throughout the mixing process there is also a direct reaction between NO and NH3 to form N2. In each step part of the fixed nitrogen (originally NO) was converted to N2. This is the goal of the reburn process.\nSometimes a NOx. reduction process is used in which the upper furnace fuel is not added in sufficient qualities to consume all of the oxygen remaining in the gas after the initial combustion. In such a process it is necessary that large volumes become reducing while parallel volumes remain oxidizing. In the reducing volumes N2, NH3, and HCN are formed. Then the reducing and oxidizing gases mix together and the remainder of the fuel is consumed. At this point the reduced nitrogen species are oxidized to N2 and NO. Again there is direct reduction to N2 by the reaction between NH3 and NO.\nThe process which is sometimes called controlled fuel lean reburn usually requires natural gas as the upper furnace fuel. Natural gas is expensive. Penetration and mixing is a great problem. Utility boiler furnaces have horizontal dimensions of 50 feet and greater. The carrier gas may be steam, air, or recycled combustion products. Often it is necessary to use a carrier gas to assure adequate penetration of the natural gas into the furnace. If the upper furnace natural gas is 5% of the fuel and the fuel is only 10% of the air flow, upper furnace injected natural gas is perhaps only 0.5% of the gas flow. The combustion products being quite hot may have a volume as high as 1000 times the upper furnace natural gas. Use of steam as a carrier gas is expensive. The use of air or recycled combustion products requires expensive duct work. Often there is no place for the duct work. The boiler face is simply too crowded with necessary equipment to allow the duct work to be installed. Large penetrations through the furnace walls are needed and this requires bending water wall tubes. The flue gas needs to be returned from a remote part of the boiler. Fans are needed for flue gas and often for air. Because air has oxygen in it, use of air as the carrier gas requires more upper furnace fuel before the gas stream can be made reducing.\nSome operators have tried coal as a reburn fuel. The burnout times for coal are longer than for natural gas. This requires that both the fuel and the burnout air be added sooner. As a result, much of the reaction occurs at higher temperatures which results in more NOx emissions. The use of coal requires that there be additional pipes to carry primary air and pulverized fuel from the mills usually at ground level to the height where the reburn fuel is injected. It may even require an additional pulverizer.\nA method of reducing NOx by injecting: a coal water slurry as a fuel lean reburn fuel has been patented, as U.S. Pat. No. 5,746,144, and that invention overcomes many of the objections to reburning with coal. Yet, this fuel requires longer burnout times than natural gas. Hence, there is still a need for a reburn fuel which has the benefits of a coal water slurry while also having a shorter burnout time.\nWe provide a method of reducing NOx emissions by injecting a biowaste water slurry, a biomass water slurry, waste rubber water slurry, waste plastic water slurry, orremulsion, or a wood water slurry into the upper furnace forming a fuel lean reburn process. Preferably the slurry is injected into a region of the furnace which is 1800xc2x0 F. to 2700xc2x0 F.\nWe also prefer to add sufficient water in the slurry to drive the reaction between water and carbon monoxide to the products to the hydrogen and carbon dioxide.\nWe also prefer to add a calcium compound to the fuel slurry. The calcium reacts with the sulfur dioxide to form calcium sulfate and thus reduce the emissions of sulfur dioxide.\nWe further prefer to introduce our slurry into the furnace through a combination of atomizing nozzles and jets.\nWe may also add ammonia, urea or other fixed nitrogen compound to the fuel water slurry as a selective reducing agent to reduce NO.\nWe may add animal wastes to increase the nitrogen in the reburn slurry.\nWe may also use a catalyst to improve the reaction between the reduced fixed nitrogen and nitrogen oxides.\nOther objects and advantages of the invention will become apparent from a description of certain preferred embodiments described with reference to the figures."}
{"text": "Some display devices create an image by scanning a beam of varying intensity across a display surface. For example, in cathode ray tubes (CRTs), an electron beam is scanned across a surface in a row and column pattern. Further, some projection display devices scan a light beam across a surface in a row and column pattern. In these display devices, the beam intersects each pixel location as it paints the image row by row. The intensity of the scanned light beam is then modulated as it passes over each display pixel location.\nPixels are typically displayed at times specified by a “pixel clock”. The pixel clock is typically generated by phase locking the output of a voltage controlled oscillator (VCO) to a sync signal that describes the position of the scanned beam. VCOs used for this purpose are typically sensitive analog circuits that are susceptible to electromagnetic interference (EMI)."}
{"text": "Devices for therapeutic uses are already known, in which a current of piezoelectrical origin is produced from a generator operated by hand. Such devices are used particularly for treatment of muscular pains, the current being led to flow between electrodes which are part of the device."}
{"text": "The present invention is directed to a process making it possible to replace the three chlorine atoms of cyanuric chloride by three identical substituents in a very simple and effective manner."}
{"text": "In methods or processes with respect to hearing device adjustments, the workflow of an adjusting operation i.e. sequence of adjusting steps, the adjusting steps themselves and the timing of performing these steps, are not stored. A particular adjusting or fitting operation of a hearing device is performed independent of previous adjustments or fittings of the hearing device. Thereby, most generically, the expertise from previous adjustments is lost after the adjustment is performed, particularly if subsequent adjustments are performed by different individuals, i.e. experts, at different locations. Although, adjustments and fittings of a hearing device are often associated with fitting mechanical characteristics of the respective hearing devices e.g. shape of the shell, surface characteristics of the shell etc., they may also concern adjusting signal-processing of the respective hearing device."}
{"text": "Data compression occurs in a number of contexts. It is very commonly used in communications and computer networking to store, transmit, and reproduce information efficiently. It finds particular application in the encoding of images, audio and video. Video presents a significant challenge to data compression because of the large amount of data required for each video frame and the speed with which encoding and decoding often needs to occur. The current state-of-the-art for video encoding is the ITU-T H.264/AVC video coding standard. It defines a number of different profiles for different applications, including the Main profile, Baseline profile and others. A next-generation video encoding standard is currently under development through a joint initiative of MPEG-ITU termed High Efficiency Video Coding (HEVC). The initiative may eventually result in a video-coding standard that will form part of a suite of standards referred to as MPEG-H.\nThere are a number of standards for encoding/decoding images and videos, including H.264, that use block-based coding processes. In these processes, the image or frame is divided into blocks, typically 4×4 or 8×8, and the blocks are spectrally transformed into coefficients, quantized, and entropy encoded. In many cases, the data being transformed is not the actual pixel data, but is residual data following a prediction operation. Predictions can be intra-frame, i.e. block-to-block within the frame/image, or inter-frame, i.e. between frames (also called motion prediction). It is expected that HEVC will also have these features.\nWhen spectrally transforming residual data, many of these standards prescribe the use of a discrete cosine transform (DCT) or some variant thereon. The resulting DCT coefficients are then quantized using a quantizer to produce quantized transform domain coefficients, or indices.\nThe block or matrix of quantized transform domain coefficients (sometimes referred to as a “transform unit”) is then entropy encoded using a particular context model. In H.264/AVC and in the current development work for HEVC, the quantized transform coefficients are encoded by (a) encoding a last significant coefficient position indicating the location of the last non-zero coefficient in the transform unit, (b) encoding a significance map indicating the positions in the transform unit (other than the last significant coefficient position) that contain non-zero coefficients, (c) encoding the magnitudes of the non-zero coefficients, and (d) encoding the signs of the non-zero coefficients. This encoding of the quantized transform coefficients often occupies 30-80% of the encoded data in the bitstream.\nSimilar reference numerals may have been used in different figures to denote similar components."}
{"text": "This invention relates to newly identified polypeptides and polynucleotides encoding such polypeptides, to their use in therapy and in identifying compounds which may be agonists, antagonists and/or inhibitors which are potentially useful in therapy, and to production of such polypeptides and polynucleotides.\nThe drug discovery process is currently undergoing a fundamental revolution as it embraces xe2x80x98functional genomicsxe2x80x99, that is, high throughput genome- or gene-based biology. This approach is rapidly superseding earlier approaches based on xe2x80x98positional cloningxe2x80x99. A phenotype, that is a biological function or genetic disease, would be identified and this would then be tracked back to the responsible gene, based on its genetic map position.\nFunctional genomics relies heavily on the various tools of bioinformatics to identify gene sequences of potential interest from the many molecular biology databases now available. There is a continuing need to identify and characterise further genes and their related polypeptides/proteins, as targets for drug discovery.\nIt is well established that many medically significant biological processes are mediated by proteins participating in signal transduction pathways that involve G-proteins and/or second messengers. e.g., cAMP (Lefkowitz, Nature, 1991, 351:353-354). Herein these proteins are referred to as proteins participating in pathways with G-proteins or PPG proteins. Some examples of these proteins include the GPC receptors, such as those for adrenergic agents and dopamine (Kobilka, B. K., et al., Proc. Natl Acad. Sci., USA. 1987, 84:46-50; Kobilka, B. K., et al., Science, 1987, 238:650-656; Bunzow, J. R., et al., Nature, 1988, 336:783-787), G-proteins themselves, effector proteins, e.g., phospholipase C, adenyl cyclase, and phosphodiesterase, and actuator proteins, e.g., protein kinase A and protein kinase C (Simon, M. I., et al., Science, 1991, 252:802-8).\nFor example, in one form of signal transduction, the effect of hormone binding is activation of the enzyme, adenylate cyclase, inside the cell. Enzyme activation by hormones is dependent on the presence of the nucleotide GPT. GPT also influences hormone binding. A G-protein connects the hormone receptor to adenylate cyclase. G-protein was shown to exchange GTP for bound GDP when activated by a hormone receptor. The GTP-carrying form then binds to activated adenylate cyclase. Hydrolysis of GTP to GDP, catalyzed by the G-protein itself, returns the G-protein to its basal, inactive form. Thus, the G-protein serves a dual role, as an intermediate that relays the signal from receptor to effector, and as a clock that controls the duration of the signal.\nThe membrane protein gene superfamily of G-protein coupled receptors has been characterized as having seven putative transmembrane domains. The domains are believed to represent transmembrane a-helices connected by extracellular or cytoplasmic loops. G-protein coupled receptors include a wide range of biologically active receptors, such as hormone, viral, growth factor and neuroreceptors.\nG-protein coupled receptors (otherwise known as 7TM receptors) have been characterized as including these seven conserved hydrophobic stretches of about 20 to 30 amino acids, connecting at least eight divergent hydrophilic loops. The G-protein family of coupled receptors includes dopamine receptors which bind to neuroleptic drugs used for treating psychotic and neurological disorders. Other examples of members of this family include, but are not limited to, calcitonin, adrenergic, endothelin, cAMP, adenosine, muscarinic, acetylcholine, serotonin, histamine, thrombin, kinin, follicle stimulating hormone, opsins, endothelial differentiation gene-1, rhodopsins, odorant, and cytomegalovirus receptors.\nMost G-protein coupled receptors have single conserved cysteine residues in each of the first two extracellular loops which form disulfide bonds that are believed to stabilize functional protein structure. The 7 transmembrane regions are designated as TM1, TM2, TM3, TM4, TM5, TM6, and TM7. TM3 has been implicated in signal transduction. Phosphorylation and lipidation (palmitylation or farnesylation) of cysteine residues can influence signal transduction of some G-protein coupled receptors. Most G-protein coupled receptors contain potential phosphorylation sites within the third cytoplasmic loop and/or the carboxy terminus. For several G-protein coupled receptors, such as the b-adrenoreceptor, phosphorylation by protein kinase A and/or specific receptor kinases mediates receptor desensitization. For some receptors, the ligand binding sites of G-protein coupled receptors are believed to comprise hydrophilic sockets formed by several G-protein coupled receptor transmembrane domains, said socket being surrounded by hydrophobic residues of the G-protein coupled receptors. The hydrophilic side of each G-protein coupled receptor transmembrane helix is postulated to face inward and form polar ligand binding site. TM3 has been implicated in several G-protein coupled receptors as having a ligand binding site, such as the TM3 aspartate residue. TM5 serines, a TM6 asparagine and TM6 or TM7 phenylalanines or tyrosines are also inmplicated in ligand binding.\nG-protein coupled receptors can be intracellulary coupled by heterotrimeric G-proteins to various intracellular enzymes, ion channels and transporters (see, Johnson et al., Endoc. Rev., 1989, 10:317-331) Different G-protein a-subunits preferentially stimulate particular effectors to modulate various biological functions in a cell. Phosphorylation of cytoplasmic residues of G-protein coupled receptors have been identified as an important mechanism for the regulation of G-protein couplings of some G-protein coupled receptors. G-protein coupled receptors are found in numerous sites within a mammalian host. Over the past 15 years, nearly 350 therapeutic agents targeting 7 transmembrane (7 TM) receptors have been successfully introduced onto the market.\nThe present invention relates to HOFNH30, in particular HOFNH30 polypeptides and HOFNH30 polynucleotides, recombinant materials and methods for their production. In another aspect, the invention relates to methods for using such polypeptides and polynucleotides, including the treatment of infections such as bacterial, fungal, protozoan and viral infections, particularly infections caused by HIV-1 or HIV-2; cancers; diabetes; asthma; Parkinson\"\"s disease; both acute and congestive heart failure; hypotension; hypertension; urinary retention; osteoporosis; angina pectoris; myocardial infarction; ulcers; asthma; allergies; benign prostatic hypertrophy (benign prostatic hyperplasia); chronic renal failure; renal disease; impaired glucose tolerance; seizure disorder; depression; anxiety; obsessive compulsive disorder; affective neurosis/disorder; depressive neurosis/disorder; anxiety neurosis; dysthymic disorder; behavior disorder; mood disorder; shizophrenia; psychosexual dysfunction; sex disorder; sexual disorder; disturbed biological and circadian rhythms; feeding disorders, such as anorexia, bulimia, cachexia, and obesity; Cushing\"\"s syndrome/disease; basophil adenoma; prolactinoma; hyperprolactinemia; hypopituitarism; hypophysis tumor/adenoma; hypothalamic diseases; Froehlich\"\"s syndrome; adenohypophysis disease; hypophysis disease; hypophysis tumor/adenoma; pituitary growth hormone; adenohypophysis hypofunction; adrenohpophysis hyperfunction; hypothalamic hypogonadism; Kallman\"\"s syndrome (anosmia, hyposmia); functional or psychogenic amenorrhea; hypopituitarism; hypothalamic hypothyroidism; hypothalamic-adrenal dysfunction; idiopathic hyperprolactinemia; hypothalamic disorders of growth hormone deficiency; idiopathic growth hormone deficiency; dwarfism; gigantism; acromegaly; disturbed biological and circadian rhythms; and sleep disturbances associated with such diseases as neurological disorders, heart and lung diseases, mental illness, and addictions: migraine; hyperalgesia; enhanced or exaggerated sensitivity to pain, such as hyperalgesia, causalgia and allodynia; acute pain; burn pain; atypical facial pain; neuropathic pain; back pain; complex regional pain syndromes I and II; arthritic pain; sports injury pain; pain related to infection, e.g., HIV, post-polio syndrome, and post-herpetic neuralgia; phantom limb pain; labour pain; cancer pain; post-chemotherapy pain; post-stroke pain; post-operative pain; neuralgia; and tolerance to narcotics or withdrawal from narcotics; sleep disorders; sleep apnea; narcolepsy; insomnia; parasomnia; jet-lag syndrome; and other neurodegenerative disorders, which includes nosological entities such as disinhibition-dementia-parkinsonism-amyotrophy complex; pallido-ponto-nigral degeneration; and dyskinesias, such as Huntington\"\"s disease or Gilles dela Tourett\"\"s syndrome, hereinafter referred to as xe2x80x9cthe Diseasesxe2x80x9d, amongst others. In a still fturther aspect, the invention relates to diagnostic assays for detecting Diseases associated with inappropriate HOFNH30 activity or levels.\nIn accordance with another aspect of the present invention there are provided methods of screening for compounds which bind to and activate the HOFNH30 polypeptides (receptors) of the present invention (called agonists), or inhibit the interaction of the HOFNH30 polypeptides with receptor ligands (called antagonists). In particular, a preferred method for identifying agonist or antagonist of a receptor of the present invention comprises:\n(a) contacting a cell expressing on the surface thereof the receptor, said receptor being associated with a second component capable of providing a detectable signal in response to the binding of a compound to said receptor, with a compound to be screened under conditions to permit binding to the receptor; and\n(b) determining whether the compound binds to and activates or inhibits the receptor by measuring the level of a signal generated from the interaction of the compound with the receptor.\nFurther, preferred method for identifying agonist or antagonist of a receptor of the present invention comprises:\n(a) contacting a cell expressing on the surface thereof the receptor, said receptor being associated with a second component capable of providing a detectable signal in response to the binding of a compound to said receptor, with a compound to be screened under conditions to permit binding to the receptor; and\n(b) determining whether the compound binds to and activates or inhibits the receptor by comparing the level of a signal generated from the interaction of the compound with the receptor with the level of a signal without the presence of the compound.\nIn further preferred embodiments, the two methods described above further comprises conducting the identification of agonist or antagonist in the presence of labeled or unlabeled LPA.\nIn another embodiment of the method for identifying agonist or antagonist of a receptor of the present invention comprises:\ndetermining the inhibition of binding of a ligand to cells which have the receptor on the surface thereof, or to cell membranes containing the receptor, in the presence of a candidate compound under conditions to permit binding to the receptor, and determining the amount of ligand bound to the receptor, such that a compound capable of causing reduction of binding of a ligand is an agonist or antagonist. Preferably the ligand is LPA. Yet more preferably LPA is labeled. More particularly, a method of screening for HOFNH30 receptor antagonist or agonist comprises the steps of:\n(a) incubating a labeled LPA with a whole cell expressing HOFNH30 receptor on the cell surface, or cell membrane containing HOFNH30 receptor;\n(b) measuring the amount of labeled LPA bound to the whole cell or the cell membrane;\n(c) adding a candidate compound to a mixture of labeled LPA and the whole cell or the cell membrane of step (a) and allowing to attain equilibrium;\n(d) measuring the amount of labeled LPA bound to the whole cell or the cell membrane after step (c); and\n(e) comparing the difference in the labeled LPA bound in step (b) and (d), such that the compound which causes the reduction in binding in step (d) is an agonist or antagonist.\nIn still another aspect, the invention relates to agonists and antagonists discovered by the present screening methods. Yet in still another aspect, the invention relates to treating conditions associated with HOFNH30 imbalance with the identified compounds.\nIn a first aspect, the present invention relates to HOFNH30 polypeptides. Such peptides include isolated polypeptides comprising an amino acid sequence which has at least 70% identity, preferably at least 80% identity, more preferably at least 90% identity, yet more preferably at least 95% identity, most preferably at least 97-99% identity, to that of SEQ ID NO:2 over the entire length of SEQ ID NO:2. Such polypeptides include those comprising the amino acid of SEQ ID NO:2.\nFurther peptides of the present invention include isolated polypeptides in which the amino acid sequence has at least 70% identity. preferably at least 80% identity, more preferably at least 90% identity, yet more preferably at least 95% identity, most preferably at least 97-99% identity, to the amino acid sequence of SEQ ID NO:2 over the entire length of SEQ ID NO:2. Such polypeptides include the polypeptide of SEQ ID NO:2.\nFurther peptides of the present invention include isolated polypeptides encoded by a polynucleotide comprising the sequence contained in SEQ ID NO:1.\nPolypeptides of the present invention are believed to be members of the EDG family of polypeptides. They are therefore of interest because this invention adds an additional member in the EDG family of genes that are involved in a number of biological and diseased manifestations. These properties are hereinafter referred to as xe2x80x9cHOFNH30 activityxe2x80x9d or xe2x80x9cHOFNH30 polypeptide activityxe2x80x9d or xe2x80x9cbiological activity of HOFNH30xe2x80x9d. Also included amongst these activities are antigenic and immunogenic activities of said HOFNH30 polypeptides, in particular the antigenic and immunogenic activities of the polypeptide of SEQ ID NO:2. Preferably, a polypeptide of the present invention exhibits at least one biological activity of HOFNH30.\nThe polypeptides of the present invention may be in the form of the xe2x80x9cmaturexe2x80x9d protein or may be a part of a larger protein such as a fusion protein. It is often advantageous to include an additional amino acid sequence which contains secretory or leader sequences, pro-sequences, sequences which aid in purification such as multiple histidine residues, or an additional sequence for stability during recombinant production.\nThe present invention also includes include variants of the aforementioned polypeptides, that is polypeptides that vary from the referents by conservative amino acid substitutions, whereby a residue is substituted by another with like characteristics. Typical such substitutions are among Ala, Val, Leu and Ile; among Ser and Thr; among the acidic residues Asp and Glu; among Asn and Gln; and among the basic residues Lys and Arg; or aromatic residues Phe and Tyr. Particularly preferred are variants in which several, 5-10, 1-5, 1-3, 1-2 or 1 amino acids are substituted, deleted, or added in any combination.\nPolypeptides of the present invention can be prepared in any suitable manner. Such polypeptides include isolated naturally occurring polypeptides, recombinantly produced polypeptides, synthetically produced polypeptides, or polypeptides produced by a combination of these methods. Means for preparing such polypeptides are well understood in the art.\nIn a further aspect, the present invention relates to HOFNH30 polynucleotides. Such polynucleotides include isolated polynucleotides comprising a nucleotide sequence encoding a polypeptide which has at least 70% identity, preferably at least 80% identity, more preferably at least 90% identity, yet more preferably at least 95% identity, to the amino acid sequence of SEQ ID NO:2, over the entire length of SEQ ID NO:2. In this regard, polypeptides which have at least 97% identity are highly preferred, whilst those with at least 98-99% identity are more highly preferred, and those with at least 99% identity are most highly preferred. Such polynucleotides include a polynucleotide comprising the nucleotide sequence contained in SEQ ID NO:1 encoding the polypeptide of SEQ ID NO:2.\nFurther polynucleotides of the present invention include isolated polynucleotides comprising a nucleotide sequence that has at least 70% identity, preferably at least 80% identity, more preferably at least 90% identity, yet more preferably at least 95% identity, to a nucleotide sequence encoding a polypeptide of SEQ ID NO:2, over the entire coding region. In this regard, polynucleotides which have at least 97% identity are highly preferred, whilst those with at least 98-99% identity are more highly preferred, and those with at least 99% identity are most highly preferred.\nFurther polynucleotides of the present invention include isolated polynucleotides comprising a nucleotide sequence which has at least 70% identity, preferably at least 80% identity, more preferably at least 90% identity, yet more preferably at least 95% identity, to SEQ ID NO:1 over the entire length of SEQ ID NO:1. In this regard, polynucleotides which have at least 97% identity are highly preferred, whilst those with at least 98-99% identify are more highly preferred, and those with at least 99% identity are most highly preferred. Such polynucleotides include a polynucleotide comprising the polynucleotide of SEQ ID NO:1 as well as the polynucleotide of SEQ ID NO:1.\nThe invention also provides polynucleotides which are complementary to all the above described polynucleotides.\nThe nucleotide sequence of SEQ ID NO:1 shows homology with EDG2. The nucleotide sequence of SEQ ID NO:1 is a cDNA sequence and comprises a polypeptide encoding sequence (nucleotide1 to 1065) encoding a polypeptide of 354 amino acids, the polypeptide of SEQ ID NO:2. The nucleotide sequence encoding the polypeptide of SEQ ID NO:2 may be identical to the polypeptide encoding sequence contained in SEQ ID NO:1 or it may be a sequence other than the one contained in SEQ ID NO:1, which, as a result of the redundancy (degeneracy) of the genetic code, also encodes the polypeptide of SEQ ID NO:2. The polypeptide of the SEQ ID NO:2 is structurally related to other proteins of the EDG family, having homology and/or structural similarity with EDG2. Preferred polypeptides and polynucleotides of the present invention are expected to have, inter alia, similar biological functions/properties to their homologous polypeptides and polynucleotides. Furthermore, preferred polypeptides and polynucleotides of the present invention have at least one HOFNH30 activity.\nPolynucleotides of the present invention may be obtained, using standard cloning and screening techniques, from a cDNA library derived from mRNA in cells of human Placenta, using the expressed sequence tag (EST) analysis (Adams, M. D., et al. Science (1991) 252:1651-1656; Adams, M. D. et al., Nature, (1992) 355:632-634; Adams, M. D., et al., Nature (1995) 377 Supp:3-174). Polynucleotides of the invention can also be obtained from natural sources such as genomic DNA libraries or can be synthesized using well known and commercially available techniques.\nWhen polynucleotides of the present invention are used for the recombinant production of polypeptides of the present invention, the polynucleotide may include the coding sequence for the mature polypeptide, by itself; or the coding sequence for the mature polypeptide in reading frame with other coding sequences, such as those encoding a leader or secretory sequence, a pre-, or pro- or prepro-protein sequence, or other fusion peptide portions. For example, a marker sequence which facilitates purification of the fused polypeptide can be encoded. In certain preferred embodiments of this aspect of the invention, the marker sequence is a hexa-histidine peptide, as provided in the pQE vector (Qiagen, Inc.) and described in Gentz et al., Proc Natl Acad Sci USA (1989) 86:821-824, or is an HA tag. The polynucleotide may also contain non-coding 5xe2x80x2 and 3xe2x80x2 sequences, such as transcribed, non-translated sequences, splicing and polyadenylation signals, ribosome binding sites and sequences that stabilize mRNA.\nFurther embodiments of the present invention include polynucleotides encoding polypeptide variants which comprise the amino acid sequence of SEQ ID) NO:2 and in which several, for instance from 5 to 10, 1 to 5, 1 to 3, 1 to 2 or 1, amino acid residues are substituted, deleted or added, in any combination.\nPolynucleotides which are identical or sufficiently identical to a nucleotide sequence contained in SEQ ID NO:1, may be used as hybridization probes for cDNA and genomic DNA or as primers for a nucleic acid amplification (PCR) reaction, to isolate full-length cDNAs and genomic clones encoding polypeptides of the present invention and to isolate cDNA and genomic clones of other genes (including genes encoding homologs and orthologs from species other than human) that have a high sequence similarity to SEQ ID NO:1. Typically these nucleotide sequences are 70% identical, preferably 80% identical, more preferably 90% identical, most preferably 95% identical to that of the referent. The probes or primers will generally comprise at least 15 nucleotides, preferably, at least 30 nucleotides and may have at least 50 nucleotides. Particularly preferred probes will have between 30 and 50 nucleotides.\nA polynucleotide encoding a polypeptide of the present invention, including homologs and orthologs from species other than human, may be obtained by a process which comprises the steps of screening an appropriate library under stringent hybridization conditions with a labeled probe having the sequence of SEQ ID NO: 1 or a fragment thereof; and isolating full-length cDNA and genomic clones containing said polynucleotide sequence. Such hybridization techniques are well known to the skilled artisan. Preferred stringent hybridization conditions include overnight incubation at 42xc2x0 C. in a solution comprising: 50% formamide, 5xc3x97SSC (150 mM NaCl, 15 mM trisodium citrate), 50 mM sodium phosphate (pH7.6), 5xc3x97Denhardt\"\"s solution, 10 % dextran sulfate, and 20 microgram/ml denatured, sheared salmon sperm DNA; followed by washing the filters in 0.1xc3x97SSC at about 65xc2x0 C. Thus the present invention also includes polynucleotides obtainable by screening an appropriate library under stringent hybridization conditions with a labeled probe having the sequence of SEQ ID NO:1 or a fragment thereof.\nThe skilled artisan will appreciate that, in many cases, an isolated cDNA sequence will be incomplete, in that the region coding for the polypeptide is cut short at the 5xe2x80x2 end of the cDNA. This is a consequence of reverse transcriptase, an enzyme with inherently low xe2x80x98processivityxe2x80x99 (a measure of the ability of the enzyme to remain attached to the template during the polymerisation reaction), failing to complete a DNA copy of the mRNA template during 1st strand cDNA synthesis.\nThere are several methods available and well known to those skilled in the art to obtain full-length cDNAs, or extend short cDNAs, for example those based on the method of Rapid Amplification of cDNA ends (RACE) (see, for example, Frohman et al., PNAS USA 85, 8998-9002, 1988). Recent modifications of the technique, exemplified by the Marathon(trademark) technology (Clontech Laboratories Inc.) for example, have significantly simplified the search for longer cDNAs. In the Marathon(trademark) technology, cDNAs have been prepared from mRNA extracted from a chosen tissue and an xe2x80x98adaptorxe2x80x99 sequence ligated onto each end. Nucleic acid amplification (PCR) is then carried out to amplify the xe2x80x98missingxe2x80x99 5xe2x80x2 end of the cDNA using a combination of gene specific and adaptor specific oligonucleotide primers. The PCR reaction is then repeated using xe2x80x98nestedxe2x80x99 primers, that is, primers designed to anneal within the amplified product (typically an adaptor specific primer that anneals further 3xe2x80x2 in the adaptor sequence and a gene specific primer that anneals further 5xe2x80x2 in the known gene sequence). The products of this reaction can then be analyzed by DNA sequencing and a full-length cDNA constructed either by joining the product directly to the existing cDNA to give a complete sequence, or carrying out a separate full-length PCR using the new sequence information for the design of the 5xe2x80x2 primer.\nRecombinant polypeptides of the present invention may be prepared by processes well known in the art from genetically engineered host cells comprising expression systems. Accordingly, in a further aspect, the present invention relates to expression systems which comprise a polynucleotide or polynucleotides of the present invention, to host cells which are genetically engineered with such expression systems and to the production of polypeptides of the invention by recombinant techniques. Cell-free translation systems can also be employed to produce such proteins using RNAs derived from the DNA constructs of the present invention.\nFor recombinant production, host cells can be genetically engineered to incorporate expression systems or portions thereof for polynucleotides of the present invention. Introduction of polynucleotides into host cells can be effected by methods described in many standard laboratory manuals, such as Davis et al., Basic Methods in Molecular Biology (1986) and Sambrook et al., Molecular Cloning: A Laboratory Manual, 2nd Ed., Cold Spring Harbor Laboratory Press, Cold Spring Harbor, N.Y. (1989). Preferred such methods include, for instance, calcium phosphate transfection, DEAE-dextran mediated transfection, transvection, microinjection, cationic lipid-mediated transfection, electroporation, transduction, scrape loading, ballistic introduction or infection.\nRepresentative examples of appropriate hosts include bacterial cells, such as streptococci, staphylococci, E. coli, Streptomyces and Bacillus subtilis cells; fungal cells, such as yeast cells and Aspergillus cells; insect cells such as Drosophila S2 and Spodoptera Sf9 cells; animal cells such as CHO, COS, HeLa, C127, 3T3, BHK, HEK 293 and Bowes melanoma cells; and plant cells.\nA great variety of expression systems can be used, for instance, chromosomal, episomal and virus-derived systems, e.g., vectors derived from bacterial plasmids, from bacteriophage, from transposons, from yeast episomes, from insertion elements, from yeast chromosomal elements, from viruses such as baculoviruses, papova viruses, such as SV40, vaccinia viruses, adenoviruses, fowl pox viruses, pseudorabies viruses and retroviruses, and vectors derived from combinations thereof, such as those derived from plasmid and bacteriophage genetic elements, such as cosmids and phagemids. The expression systems may contain control regions that regulate as well as engender expression. Generally, any system or vector which is able to maintain, propagate or express a polynucleotide to produce a polypeptide in a host may be used. The appropriate nucleotide sequence may be inserted into an expression system by any of a variety of well-known and routine techniques, such as, for example, those set forth in Sambrook et al., MOLECULAR CLONING, A LABORATORY MANUAL (supra). Appropriate secretion signals may be incorporated into the desired polypeptide to allow secretion of the translated protein into the lumen of the endoplasmic reticulum, the periplasmic space or the extracellular environment. These signals may be endogenous to the potypeptide or they may be heterologous signals.\nIf a polypeptide of the present invention is to be expressed for use in screening assays, it is generally preferred that the polypeptide be produced at the surface of the cell. In this event, the cells may be harvested prior to use in the screening assay. If the polypeptide is secreted into the medium, the medium can be recovered in order to recover and purify the polypeptide. If produced intracellularly, the cells must first be lysed before the polypeptide is recovered.\nPolypeptides of the present invention can be recovered and purified from recombinant cell cultures by well-known methods including ammonium sulfate or ethanol precipitation, acid extraction, anion or cation exchange chromatography, phosphocellulose chromatography, hydrophobic interaction chromatography, affinity chromatography, hydroxylapatite chromatography and lectin chromatography. Most preferably, high performance liquid chromatography is employed for purification. Well known techniques for refolding proteins may be employed to regenerate active conformation when the polypeptide is denatured during isolation and or purification.\nThis invention also relates to the use of polynucleotides of the present invention as diagnostic reagents. Detection of a mutated form of the gene characterized by the polynucleotide of SEQ ID NO:1 which is associated with a dysfunction will provide a diagnostic tool that can add to, or define, a diagnosis of a disease, or susceptibility to a disease, which results from under-expression, over-expression or altered expression of the gene. Individuals carrying mutations in the gene may be detected at the DNA level by a variety of techniques.\nNucleic acids for diagnosis may be obtained from a subject\"\"s cells, such as from blood, urine, saliva, tissue biopsy or autopsy material. The genomic DNA may be used directly for detection or may be amplified enzymatically by using PCR or other amplification techniques prior to analysis. RNA or cDNA may also be used in similar fashion. Deletions and insertions can be detected by a change in size of the amplified product in comparison to the normal genotype. Point mutations can be identified by hybridizing amplified DNA to labeled HOFNH30 nucleotide sequences. Perfectly matched sequences can be distinguished from mismatched duplexes by RNase digestion or by differences in melting temperatures. DNA sequence differences may also be detected by alterations in electrophoretic mobility of DNA fragments in gels, with or without denaturing agents, or by direct DNA sequencing (e.g., Myers et al., Science (1985) 230:1242). Sequence changes at specific locations may also be revealed by nuclease protection assays, such as RNase and S1 protection or the chemical cleavage method (see Cotton et al., Proc Natl Acad Sci USA (1985) 85: 4397-4401). In another embodiment, an array of oligonucleotides probes comprising HOFNH30 nucleotide sequence or fragments thereof can be constructed to conduct efficient screening of e.g., genetic mutations. Array technology methods are well known and have general applicability and can be used to address a variety of questions in molecular genetics including gene expression, genetic linkage, and genetic variability (see for example: M. Chee et al., Science, Vol 274, pp 610-613 (1996)).\nThe diagnostic assays offer a process for diagnosing or determining a susceptibility to the Diseases through detection of mutation in the HOFNH30 gene by the methods described. In addition, such diseases may be diagnosed by methods comprising determining from a sample derived from a subject an abnormally decreased or increased level of polypeptide or mRNA. Decreased or increased expression can be measured at the RNA level using any of the methods well known in the art for the quantitation of polynucleotides, such as, for example, nucleic acid amplification, for instance PCR, RT-PCR, RNase protection, Northern blotting and other hybridization methods. Assay techniques that can be used to determine levels of a protein, such as a polypeptide of the present invention, in a sample derived from a host are well-known to those of skill in the art. Such assay methods include radioimmunoassays, competitive-binding assays, Western Blot analysis and ELISA assays.\nThus in another aspect, the present invention relates to a diagnostic kit which comprises:\n(a) a polynucleotide of the present invention, preferably the nucleotide sequence of SEQ ID NO: 1, or a fragment thereof;\n(b) a nucleotide sequence complementary to that of (a);\n(c) a polypeptide of the present invention, preferably the polypeptide of SEQ ID NO:2 or a fragment thereof; or\n(d) an antibody to a polypeptide of the present invention, preferably to the polypeptide of SEQ ID NO:2.\nIt will be appreciated that in any such kit, (a), (b), (c) or (d) may comprise a substantial component. Such a kit will be of use in diagnosing a disease or susceptibility to a disease, particularly infections such as bacterial, fungal, protozoan and viral infections, particularly infections caused by HIV-1 or HIV-2; cancers; diabetes; asthma; Parkinson\"\"s disease; both acute and congestive heart failure; hypotension; hypertension; urinary retention; osteoporosis; angina pectoris; myocardial infarction; ulcers; asthma; allergies; benign prostatic hypertrophy (benign prostatic hyperplasia); chronic renal failure; renal disease; impaired glucose tolerance; seizure disorder; depression; anxiety; obsessive compulsive disorder; affective neurosis/disorder; depressive neurosis/disorder; anxiety neurosis; dysthymic disorder; behavior disorder; mood disorder; shizophrenia; psychosexual dysfunction; sex disorder; sexual disorder; disturbed biological and circadian rhythms; feeding disorders, such as anorexia,. bulimia, cachexia, and obesity; Cushing\"\"s syndrome/disease; basophil adenoma; prolactinoma; hyperprolactinemia; hypopituitarism; hypophysis tumor/adenoma; hypothalamic diseases; Froehlich\"\"s syndrome; adenohypophysis disease; hypophysis disease; hypophysis tumor/adenoma; pituitary growth hormone; adenohypophysis hypofunction; adrenohpophysis hyperfunction; hypothalamic hypogonadism; Kallman\"\"s syndrome (anosmia, hyposmia); functional or psychogenic amenorrhea; hypopituitarism; hypothalamic hypothyroidism; hypothalamic-adrenal dysfunction; idiopathic hyperprolactinemia; hypothalamic disorders of growth hormone deficiency; idiopathic growth hormone deficiency; dwarfism; gigantism; acromegaly; disturbed biological and circadian rhythms; and sleep disturbances associated with such diseases as neurological disorders, heart and lung diseases, mental illness, and addictions; migraine; hyperalgesia; enhanced or exaggerated sensitivity to pain, such as hyperalgesia, causalgia and allodynia; acute pain; burn pain; atypical facial pain; neuropathic pain; back pain; complex regional pain syndromes I and II; arthritic pain; sports injury pain; pain related to infection, e.g., HIV, post-polio syndrome, and post-herpetic neuralgia; phantom limb pain; labour pain; cancer pain; post-chemotherapy pain; post-stroke pain; post-operative pain; neuralgia; and tolerance to narcotics or withdrawal from narcotics; sleep disorders; sleep apnea; narcolepsy; insomnia; parasomrnia; jet-lag syndrome; and other neurodegenerative disorders, which includes nosological entities such as disinhibition-dementia-parkinsonism-amyotrophy complex; pallido-ponto-nigral degeneration; and dyskinesias, such as Huntington\"\"s disease or Gilles dela Tourett\"\"s syndrome, amongst others.\nThe nucleotide sequences of the present invention are also valuable for chromosome identification. The sequence is specifically targeted to, and can hybridize with, a particular location on an individual human chromosome. The mapping of relevant sequences to chromosomes according to the present invention is an important first step in correlating those sequences with gene associated disease. Once a sequence has been mapped to a precise chromosomal location, the physical position of the sequence on the chromosome can be correlated with genetic map data. Such data are found in, for example, V. McKusick, Mendelian Inheritance in Man (available on-line through Johns Hopkins University Welch Medical Library). The relationship between genes and diseases that have been mapped to the same chromosomal region are then identified through linkage analysis (coinheritance of physically adjacent genes).\nThe differences in the cDNA or genomic sequence between affected and unaffected individuals can also be determined. If a mutation is observed in some or all of the affected individuals but not in any normal individuals, then the mutation is likely to be the causative agent of the disease.\nThe polypeptides of the invention or their fragments or analogs thereof, or cells expressing them, can also be used as immunogens to produce antibodies immunospecific for polypeptides of the present invention. The term xe2x80x9cimmunospecificxe2x80x9d means that the antibodies have substantially greater affinity for the polypeptides of the invention than their affinity for other related polypeptides in the prior art.\nAntibodies generated against polypeptides of the present invention may be obtained by administering the polypeptides or epitope-bearing fragments, analogs or cells to an animal, preferably a non-human animal, using routine protocols. For preparation of monoclonal antibodies, any technique which provides antibodies produced by continuous cell line cultures can be used. Examples include the hybridoma technique (Kohler, G. and Milstein, C., Nature (1975) 256:495-497), the trioma technique, the human B-cell hybridoma technique (Kozbor et al., Immunology Today (1983) 4:72) and the EBV-hybridoma technique (Cole et al., MONOCLONAL ANTIBODIES AND CANCER THERAPY, pp. 77-96, Alan R. Liss, Inc., 1985).\nTechniques for the production of single chain antibodies, such as those described in U.S. Pat. No. 4,946,778, can also be adapted to produce single chain antibodies to polypeptides of this invention. Also, transgenic mice, or other organisms, including other mammals, may be used to express humanized antibodies.\nThe above-described antibodies may be employed to isolate or to identify clones expressing the polypeptide or to purify the polypeptides by affinity chromatography.\nAntibodies against polypeptides of the present invention may also be employed to treat the Diseases, amongst others.\nIn a further aspect, the present invention relates to genetically engineered soluble fusion proteins comprising a polypeptide of the present invention, or a fragment thereof, and various portions of the constant regions of heavy or light chains of immunoglobulins of various subclasses (IgG, IgM, IgA, IgE). Preferred as an immunoglobulin is the constant part of the heavy chain of human IgG, particularly IgG1, where fusion takes place at the hinge region. In a particular embodiment, the Fc part can be removed simply by incorporation of a cleavage sequence which can be cleaved with blood clotting factor Xa. Furthermore, this invention relates to processes for the preparation of these fusion proteins by genetic engineering, and to the use thereof for drug screening, diagnosis and therapy. A further aspect of the invention also relates to polynucleotides encoding such fusion proteins. Examples of fusion protein technology can be found in International Patent Application Nos. WO94/29458 and WO94/22914.\nAnother aspect of the invention relates to a method for inducing an immunological response in a mammal which comprises inoculating the mammal with a polypeptide of the present invention, adequate to produce antibody and/or T cell immune response to protect said animal from the Diseases hereinbefore mentioned, amongst others. Yet another aspect of the invention relates to a method of inducing immunological response in a mammal which comprises, delivering a polypeptide of the present invention via a vector directing expression of the polynucleotide and coding for the polypeptide in vivo in order to induce such an immunological response to produce antibody to protect said animal from diseases.\nA further aspect of the invention relates to an immunological/vaccine formulation (composition) which, when introduced into a mammalian host, induces an immunological response in that mammal to a polypeptide of the present invention wherein the composition comprises a polypeptide or polynucleotide of the present invention. The vaccine formulation may further comprise a suitable carrier. Since a polypeptide may be broken down in the stomach, it is preferably administered parenterally (for instance, subcutaneous, intramuscular, intravenous, or intradermal injection). Formulations suitable for parenteral administration include aqueous and non-aqueous sterile injection solutions which may contain anti-oxidants, buffers, bacteriostats and solutes which render the formulation instonic with the blood of the recipient; and aqueous and non-aqueous sterile suspensions which may include suspending agents or thickening agents. The formulations may be presented in unit-dose or multi-dose containers, for example, sealed ampoules and vials and may be stored in a freeze-dried condition requiring only the addition of the sterile liquid carrier immediately prior to use. The vaccine formulation may also include adjuvant systems for enhancing the immunogenicity of the formulation, such as oil-in water systems and other systems known in the art. The dosage will depend on the specific activity of the vaccine and can be readily determined by routine experimentation.\nWe have now discovered that LPA is a ligand for HOFNH30 polypeptide. The HOFNH30 polypeptide of the present invention may be employed in a process for screening for compounds which bind to and activate the HOFNH30 polypeptides of the present invention (called agonists), or inhibit the interaction of the HOFNH30 polypeptides with receptor ligands (called antagonists).\nThus, polypeptides of the invention may also be used to assess the binding of small molecule substrates and ligands in, for example, cells, cell-free preparations, chemical libraries, and natural product mixtures. These substrates and ligands may be natural substrates and ligands or may be structural or functional mimetics. See Coligan et al., Current Protocols in Immunology 1(2):Chapter 5 (1991).\nHOFNH30 proteins are responsible for many biological functions, including many pathologies. Accordingly, it is desirous to find compounds and drugs which stimulate HOFNH30 on the one hand and which can inhibit the function of HOFNH30 on the other hand. In general, agonists or antagoinsts, as the case may be, are employed for therapeutic and prophylactic purposes for such conditions as infections such as bacterial, fungal, protozoan and viral infections, particularly infections caused by HIV-1 or HIV-2; cancers; diabetes; asthma; Parkinson\"\"s disease; both acute and congestive heart failure; hypotension; hypertension; urinary retention; osteoporosis; angina pectoris; myocardial infarction; ulcers; asthma; allergies; benign prostatic hypertrophy (benign prostatic hyperplasia); chronic renal failure; renal disease; impaired glucose tolerance; seizure disorder; depression; anxiety; obsessive compulsive disorder; affective neurosis/disorder; depressive neurosis/disorder; anxiety neurosis; dysthymic disorder; behavior disorder; mood disorder; shizophrenia; psychosexual dysfunction; sex disorder; sexual disorder; disturbed biological and circadian rhythms; feeding disorders; such as anorexia, bulimia, cachexia, and obesity; Cushing\"\"s syndrome/disease; basophil adenoma; prolactinoma; hyperprolactinemia; hypopituitarism; hypophysis tumor/adenoma; hypothalamic diseases; Froehlich\"\"s syndrome; adenohypophysis disease; hypophysis disease; hypophysis tumor/adenoma; pituitary growth hormone; adenohypophysis hypofunction; adrenohpophysis hyperfunction; hypothalamic hypogonadism; Kallman\"\"s syndrome (anosmia, hyposmia); functional or psychogenic amenorrhea; hypopituitarism; hypothalamic hypothyroidism; hypothalamic-adrenal dysfunction; idiopathic hyperprolactinemia; hypothalamic disorders of growth hormone deficiency; idiopathic growth hormone deficiency; dwarfism, gigantism; acromegaly; disturbed biological and circadian rhythms; and sleep disturbances associated with such diseases as neurological disorders, heart and lung diseases, mental illness, and addictions, migraine: hyperalgesia; enhanced or exaggerated sensitivity to pain, such as hyperalgesia, causalgia and allodynia; acute pain; burn pain; atypical facial pain; neuropathic pain; back pain; complex regional pain syndromes I and II; arthritic pain; sports injury pain; pain related to infection, e.g., HIV, post-polio syndrome, and post-herpetic neuralgia; phantom limb pain; labour pain; cancer pain; post-chemotherapy pain; post-stroke pain; post-operative pain; neuralgia; and tolerance to narcotics or withdrawal from narcotics; sleep disorders; sleep apnea; narcolepsy; insomnia; parasomnia; jet-lag syndrome; and other neurodegenerative disorders, which includes nosological entities such as disinhibition-dementia-parkinsonism-amyotrophy complex; pallido-ponto-nigral degeneration; and dyskinesias, such as Huntington\"\"s disease or Gilles dela Tourett\"\"s syndrome.\nIn general, such screening procedures involve providing appropriate cells which express the receptor polypeptide of the present invention on the surface thereof. Such cells include cells from mammals, yeast, Drosophila or E. coli. In particular, a polynucleotide encoding the receptor of the present invention is employed to transfect cells to thereby express the HOFNH30 polypeptide. The expressed receptor is then contacted with a test compound to observe binding, stimulation or inhibition of a functional response.\nOne such screening procedure involves the use of melanophores which are transfected to express the HOFNH30 polypeptide of the present invention. Such a screening technique is described in PCT WO 92/01810, published Feb. 6, 1992. Such an assay may be employed to screen for a compound which inhibits activation of the receptor polypeptide of the present invention by contacting the melanophore cells which encode the receptor with both the receptor ligand, such as LPA, and a compound to be screened. Inhibition of the signal generated by the ligand indicates that a compound is a potential antagonist for the receptor, i.e., inhibits activation of the receptor.\nThe technique may also be employed for screening of compounds which activate the receptor by contacting such cells with compounds to be screened and determining whether such compound generates a signal, i.e., activates the receptor. Other screening techniques include the use of cells which express the HOFNH30 polypeptide (for example, transfected CHO cells) in a system which measures extracellular pH changes caused by receptor activation. In this technique, compounds may be contacted with cells expressing the receptor polypeptide of the present invention. A second messenger response, e.g., signal transduction or pH changes, is then measured to determine whether the potential compound activates or inhibits the receptor.\nAnother screening technique involves expressing the HOFNH30 polypeptide in which the receptor is linked to phospholipase C or D. Representative examples of such cells include, but are not limited to, endothelial cells, smooth muscle cells, and embryonic kidney cells. The screening may be accomplished as hereinabove described by detecting activation of the receptor or inhibition of activation of the receptor from the phospholipase second signal.\nAnother method involves screening for compounds which are antagonists or agonists by determining inhibition of binding of labeled ligand, such as LPA, to cells which have the receptor on the surface thereof, or cell membranes containing the receptor. Such a method involves transfecting a cell (such as eukaryotic cell) with DNA encoding the HOFNH30 polypeptide such that the cell expresses the receptor on its surface. The cell is then contacted with a potential antagonist or agonist in the presence of a labeled form of a ligand, such as LPA. The ligand can be labeled, e.g., by radioactivity. The amount of labeled ligand bound to the receptors is measured, e.g., by measuring radioactivity associated with transfected cells or membrane from these cells. If the compound binds to the receptor, the binding of labeled ligand to the receptor is inhibited as determined by a reduction of labeled ligand which binds to the receptors. This method is called binding assay.\nAnother screening procedure involves the use of mammalian cells (CHO, HEK 293, Xenopus Oocytes, RBL-2H3, etc) which are tnansfected to express the receptor of interest. The cells are loaded with an indicator dye that produces a fluorescent signal when bound to calcium, and the cells are contacted with a test substance and a receptor agonist, such as LPA. Any change in fluorescent signal is measured over a defined period of time using, for example, a fluorescence spectrophotometer or a fluorescence imaging plate reader. A change in the fluorescence signal pattern generated by the ligand indicates that a compound is a potential antagonist or agonist for the receptor.\nAnother screening procedure involves use of mammalian cells (CHO, HEK293, Xenopus Oocytes, RBL-2H3, etc.) which are transfected to express the receptor of interest, and which are also transfected with a reporter gene construct that is coupled to activation of the receptor (for example, luciferase or beta-galactosidase behind an appropriate promoter). The cells are contacted with a test substance and the receptor agonist (ligand), such as LPA, and the signal produced by the reporter gene is measured after a defined period of time. The signal can be measured using a luminometer, spectrophotometer, fluorimeter, or other such instrument appropriate for the specific reporter construct used. Change of the signal generated by the ligand indicates that a compound is a potential antagonist or agonist for the receptor.\nAnother screening technique for antagonists or agonits involves introducing RNA encoding the HOFNH30 polypeptide into Xenopus oocytes (or CHO, HEK 293, RBL-2H3, etc.) to transiently or stably express the receptor. The receptor oocytes are then contacted with the receptor ligand, such as LPA, and a compound to be screened. Inhibition or activation of the receptor is then determined by detection of a signal, such as, cAMP, calcium, proton, or other ions.\nAnother method involves screening for HOFNH30 polypeptide inhibitors by determining inhibition or stimulation of HOFNH30 polypeptide-mediated cAMP and/or adenylate cyclase accumulation or dimunition. Such a method involves transiently or stably transfecting a eukaryotic cell with HOFNH30 polypeptide receptor to express the receptor on the cell surface. The cell is then exposed to potential antagonists or agonists in the presence of HOFNH30 polypeptide ligand, such as LPA. The changes in levels of cAMP is then measured over a defined period of time, for example, by radio-immuno or protein binding assays (for example using Flashplates or a scintillation proximity assay). Changes in cAMP levels can also be determined by directly measuring the activity of the enzyme, adenylyl cyclase, in broken cell preparations. If the potential antagonist or agonist binds the receptor, and thus inhibits HOFNH30 polypeptide-ligand binding, the levels of HOFNH30 polypeptide-mediated cAMP, or adenylate cyclase activity, will be reduced or increased.\nOne preferred screening method involves co-transfecting HEK-293 cells with a mammalian expression plasmid encoding a G-protein coupled receptor (GPCR), such as HOFNH30, along with a mixture comprised of mammalian expression plasmids cDNAs encoding Gxcex115 (Wilkie T. M. et al Proc Natl Acad Sci USA 1991 88:10049-10053), Gxcex116 (Amatruda T. T. et al Proc Natl Acad Sci USA 1991 8:5587-5591, and three chimeric G-proteins refered to as Gqi5, Gqs5, and Gqo5 (Conklin B R et al Nature 1993 363:274-276, Conklin B. R. et al Mol Pharmacol 1996 50:885-890). Following a 24 h incubation the transfected HEK-293 cells are plated into poly-D-lysine coated 96 well black/clear plates (Becton Dickinson, Bedford, Mass.). The cells are assayed on FLIPR (Fluorescent Imaging Plate Reader, Molecular Devices, Sunnyvale, Calif.) for a calcium mobilization response following addition of test ligands. Upon identification of a ligand which stimulates calcium mobilization in HEK-293 cells expressing a given GPCR and the G-protein mixtures, subsequent experiments are performed to determine which, if any, G-protein is required for the functional response. HEK-293 cells are then transfected with the test GPCR, or co-transfected with the test GPCR and Gxcex115, Gxcex116, Gqi5, Gqs5, or Gqo5. If the GPCR requires the presence of one of the G-proteins for functional expression in HEK-293 cells, all subsequent experiments are performed with HEK-293 cell co-transfected with the GPCR and the G-protein which gives the best response. Alternatively, the receptor can be expressed in a different cell line, for example RBL-2H3, without additional G-proteins.\nAnother screening method for agonists and antagonists relies on the endogenous pheromone response pathway in the yeast, Saccharomyces cerevisiae. Heterothallic strains of yeast can exist in two mitotically stable haploid mating types, MATa and MATa. Each cell type secretes a small peptide hormone that binds to a G-protein coupled receptor on opposite mating-type cells which triggers a MAP kinase cascade leading to GI arrest as a prelude to cell fusion. Genetic alteration of certain genes in the pheromone response pathway can alter the normal response to pheromone, and heterologous expression and coupling of human G-protein coupled receptors and humanized G-protein subunits in yeast cells devoid of endogenous pheromone receptors can be linked to downstream signaling pathways and reporter genes (e.g., U.S. Pat. Nos. 5,063,154; 5,482,835; 5,691,188). Such genetic alterations include, but are not limited to, (i) deletion of the STE2 or STE3 gene encoding the endogenous G-protein coupled pheromone receptors; (ii) deletion of the FAR1 gene encoding a protein that normally associates with cyclin-dependent kinases leading to cell cycle arrest; and (iii) construction of reporter genes fused to the FUS1 gene promoter (where FUS1 encodes a membrane-anchored glycoprotein required for cell fusion). Downstream reporter genes can permit either a positive growth selection (e.g., histidine prototrophy using the FUS1-HIS3 reporter), or a colorimetric, fluorimetric or spectrophotometnc readout, depending on the specific reporter construct used (e.g., b-galactosidase induction using a FUS1-LacZ reporter).\nThe yeast cells can be further engineered to express and secrete small peptides from random peptide libraries, some of which can permit autocrine activation of heterologously expressed human (or mammalian) G-protein coupled receptors (Broach, J. R. and Thorner, J. Nature 384: 14-16, 1996; Manfredi et al., Mol. Cell. Biol. 16: 4700-4709, 1996). This provides a rapid direct growth selection (e.g, using the FUS1-HIS3 reporter) for surrogate peptide agonists that activate characterized or orphan receptors. Alternatively, yeast cells that functionally express human (or mammalian) G-protein coupled receptors linked to a reporter gene readout (e.g., FUS1-LacZ) can be used as a platform for high-throughput screening of known ligands, fractions of biological extracts and libraries of chemical compounds for either natural or surrogate ligands. Functional agonists of sufficient potency (whether natural or surrogate) can be used as screening tools in yeast cell-based assays for identifying G-protein coupled receptor antagonists. For example, agonists will promote growth of a cell with FUS-HIS3 reporter or give positive readout for a cell with FUS1-LacZ. However, a candidate compound which inhibits growth or negates the positive readout induced by an agonist is an antagonist. For this purpose, the yeast system offers advantages over mammalian expression systems due to its ease of utility and null receptor background (lack of endogenous G-protein coupled receptors) which often interferes with the ability to identify agonists or antagonists.\nThe present invention also provides a method for identifying new ligands not known to be capable of binding to an HOFNH30 polypeptides. The screening assays described above for identifying agonists may be used to identify new ligands.\nThe present invention also contemplates agonists and antagonists obtainable from the above described screening methods.\nExamples of potential HOFNH30 polypeptide receptor antagonists include peptidomimetics, synthetic organic molecules, natural products, antibodies, etc. which bind to the receptor but do not elicit a second messenger response such that the activity of the receptor is prevented.\nPotential antagonists also include proteins which are closely related to the ligand of the HOFNH30 polypeptide receptor, i.e. a fragment of the ligand, which have lost biological function, and when they bind to the HOFNH30 polypeptide receptor, elicit no response.\nPotential antagonists also includes a small molecule which binds to the HOFNH30 polypeptide receptor, making it inaccessible to ligands such that normal biological activity is prevented. Examples of small molecules include, but are not limited to, small peptides or peptide-like molecules.\nPotential antagonists also include soluble forms of HOFNH30 polypeptide receptor, e.g., fragments of the receptor, which bind to the ligand and prevent the ligand from interacting with membrane bound HOFNH30 polypeptide receptors.\nA potential antagonist also includes an antisense construct prepared through the use of antisense technology. Antisense technology can be used to control gene expression through triple-helix formation or antisense DNA or RNA, both methods of which are based on binding of a polynucleotide to DNA or RNA. For example, the 5xe2x80x2 coding portion of the polynucleotide sequence, which encodes for the mature polypeptides of the present invention, is used to design an antisense RNA oligonucleotide of from about 10 to 40 base pairs in length. A DNA oligonucleotide is designed to be complementary to a region of the gene involved in transcription (triple helix-see Lee, et al., Nucl. Acids Res., 6: 3073 (1979); Cooney, et al, Science, 241: 456 (1988); and Dervan, et al., Science, 251: 1360 (1991)), thereby preventing transcription and production of a HOFNH30 polypeptide. The antisense RNA oligonucleotide hybridizes to the MRNA in vivo and blocks translation of the mRNA molecule to a HOFNH30 polypeptide (antisense-Okano, J., Neurochem., 56: 560 (1991); OLIGODEOXYNUCLEOTIDES AS ANTISENSE INHIBITORS OF GENE EXPRESSION, CRC Press, Boca Raton, Fla. (1988)). The oligonucleotides described above can also be delivered to cells such that the antisense RNA or DNA may be expressed in vivo to inhibit production of a HOFNH30 polypeptide.\nThus in another aspect, the present invention relates to a screening kit for identifying agonists, antagonists, and ligands for HOFNH30 polypeptides, which comprises:\n(a) a HOFNH30 polypeptide, preferably that of SEQ ID NO:2; and further preferably comprises labeled or unlabeled LPA;\n(b) a recombinant cell expressing a HOFNH30 polypeptide, preferably that of SEQ ID NO:2; and further preferably comprises labeled or unlabeled LPA; or\n(c) a cell membrane expressing HOFNH30 polypeptide; preferably that of SEQ ID NO:2; and further preferably comprises labeled or unlabeled LPA.\nIt will be appreciated that in any such kit, (a), (b), or (c) may comprise a substantial component.\nThe polynucleotides, polypeptides and antibodies to the polypeptide of the present invention may also be used to configure screening methods for detecting the effect of added compounds on the production of mRNA and polypeptide in cells. For example, an ELISA assay may be constructed for measuring secreted or cell associated levels of polypeptide using monoclonal and polyclonal antibodies by standard methods known in the art. This can be used to discover agents which may inhibit or enhance the production of polypeptide (also called antagonist or agonist, respectively) from suitably manipulated cells or tissues.\nThe polypeptide may be used to identify membrane bound or soluble receptors, if any, through standard receptor binding techniques known in the art. These include, but are not limited to, ligand binding and crosslinking assays in which the polypeptide is labeled with a radioactive isotope (for instance, 125I), chemically modified (for instance, biotinylated), or fused to a peptide sequence suitable for detection or purification, and incubated with a source of the putative receptor (cells, cell membranes, cell supernatants, tissue extracts, bodily fluids). Other methods include biophysical techniques such as surface plasmon resonance and spectroscopy. These screening methods may also be used to identify agonists and antagonists of the polypeptide which compete with the binding of the polypeptide to its receptors, if any. Standard methods for conducting such assays are well understood in the art.\nExamples of potential polypeptide antagonists include antibodies or, in some cases, oligonucleotides or proteins which are closely related to the ligands, substrates, receptors, enzymes, etc., as the case may be, of the polypeptide, e.g., a fragment of the ligands, substrates, receptors, enzymes, etc.; or small molecules which bind to the polypeptide of the present invention but do not elicit a response, so that the activity of the polypeptide is prevented.\nThus, in another aspect, the present invention relates to a screening kit for identifying agonists, antagonists, ligands, receptors, substrates, enzymes, etc. for polypeptides of the present invention; or compounds which decrease or enhance the production of such polypeptides, which comprises:\n(a) a polypeptide of the present invention;\n(b) a recombinant cell expressing a polypeptide of the present invention;\n(c) a cell membrane expressing a polypeptide of the present invention; or\n(d) antibody to a polypeptide of the present invention; which polypeptide is preferably that of SEQ ID NO:2.\nIt will be appreciated that in any such kit, (a), (b), (c) or (d) may comprise a substantial component.\nIt will be readily appreciated by the skilled artisan that a polypeptide of the present invention may also be used in a method for the structure-based design of an agonist, antagonist or inhibitor of the polypeptide, by:\n(a) determining in the first instance the three-dimensional structure of the polypeptide;\n(b) deducing the three-dimensional structure for the likely reactive or binding site(s) of an agonist, antagonist or inhibitor;\n(c) synthesizing candidate compounds that are predicted to bind to or react with the deduced binding or reactive site; and\n(d) testing whether the candidate compounds are indeed agonists, antagonists or inhibitors.\nIt will be further appreciated that this will normally be an interactive process.\nIn a further aspect, the present invention provides methods of treating abnormal conditions such as, for instance, infections such as bacterial, fungal, protozoan and viral infections, particularly infections caused by HIV-1 or HIV-2; cancers; diabetes; asthma; Parkinson\"\"s disease; both acute and congestive heart failure; hypotension; hypertension; urinary retention; osteoporosis; angina pectoris; myocardial infarction; ulcers; asthma; allergies; benign prostatic hypertrophy (benign prostatic hyperplasia); chronic renal failure; renal disease; impaired glucose tolerance; seizure disorder; depression; anxiety; obsessive compulsive disorder; affective neurosis/disorder; depressive neurosis/disorder; anxiety neurosis; dysthymic disorder; behavior disorder; mood disorder; shizophrenia; psychosexual dysfunction; sex disorder; sexual disorder; disturbed biological and circadian rhythms; feeding disorders, such as anorexia, bulimia, cachexia, and obesity; Cushing\"\"s syndrome/disease; basophil adenoma; prolactinoma; hyperprolactinemia; hypopituitarism; hypophysis tumor/adenoma; hypothalamic diseases; Froehlich\"\"s syndrome; adenohypophysis disease; hypophysis disease; hypophysis tumor/adenoma; pituitary growth hormone; adenohypophysis hypofunction; adrenohpophysis hyperfunction; hypothalamic hypogonadism; Kallman\"\"s syndrome (anosmia, hyposmia); functional or psychogenic amenorrhea; hypopituitarism; hypothalamic hypothyroidism; hypothalamic-adrenal dysfunction; idiopathic hyperprolactinemia; hypothalamic disorders of growth hormone deficiency; idiopathic growth hormone deficiency; dwarfism; gigantism; acromegaly; disturbed biological and circadian rhythms; and sleep disturbances associated with such diseases as neurological disorders, heart and lung diseases, mental illness, and addictions; migraine; hyperalgesia; enhanced or exaggerated sensitivity to pain, such as hyperalgesia, causalgia and allodynia; acute pain; burn pain; atypical facial pain; neuropathic pain; back pain; complex regional pain syndromes I and II; arthritic pain; sports injury pain; pain related to infection, e.g., HIV, post-polio syndrome, and post-herpetic neuralgia; phantom limb pain; labour pain; cancer pain; post-chemotherapy pain; post-stroke pain; post-operative pain; neuralgia; and tolerance to narcotics or withdrawal from narcotics; sleep disorders; sleep apnea; narcolepsy; insomnia; parasomnia; jet-lag syndrome; and other neurodegenerative disorders, which includes nosological entities such as disinhibition-dementia-parkinsonism-amyotrophy complex; pallido-ponto-nigral degeneration; and dyskinesias, such as Huntington\"\"s disease or Gilles dela Tourett\"\"s syndrome, related to either an excess of, or an under-expression of, HOFNH30 polypeptide activity.\nIf the activity of the polypeptide is in excess, several approaches are available. One approach comprises administering to a subject in need thereof an inhibitor compound (antagonist) as hereinabove described, optionally in combination with a pharmaceutically acceptable carrier, in an amount effective to inhibit the function of the polypeptide, such as, for example, by blocking the binding of ligands, substrates, receptors, enzymes, etc., or by inhibiting a second signal, and thereby alleviating the abnormal condition. In another approach, soluble forms of the polypeptides still capable of binding the ligand, substrate, enzymes, receptors, etc. in competition with endogenous polypeptide may be administered. Typical examples of such competitors include fragments of the HOFNH30 polypeptide.\nIn still another approach, expression of the gene encoding endogenous HOFNH30 polypeptide can be inhibited using expression blocking techniques. Known such techniques involve the use of antisense sequences, either internally generated or separately administered (see, for example, O\"\"Connor, J Neurochem (1991) 56:560 in Oligodeoxynucleotides as Antisense Inhibitors of Gene Expression. CRC Press, Boca Raton, Fla. (1988)). Alternatively, oligonucleotides which form triple helices with the gene can be supplied (see, for example, Lee et al., Nucleic Acids Res (1979) 6:3073; Cooney et al., Science (1988) 241:456; Dervan et al., Science (1991) 251:1360). These oligomers can be administered per se or the relevant oligomers can be expressed in vivo.\nFor treating abnormal conditions related to an under-expression of HOFNH30 and its activity, several approaches are also available. One approach comprises administering to a subject a therapeutically effective amount of a compound which activates a polypeptide of the present invention, i.e., an agonist as described above, in combination with a pharmaceutically acceptable carrier, to thereby alleviate the abnormal condition. Alternatively, gene therapy may be employed to effect the endogenous production of HOFNH30 by the relevant cells in the subject. For example, a polynucleotide of the invention may be engineered for expression in a replication defective retroviral vector, as discussed above. The retroviral expression construct may then be isolated and introduced into a packaging cell transduced with a retroviral plasmid vector containing RNA encoding a polypeptide of the present invention such that the packaging cell now produces infectious viral particles containing the gene of interest. These producer cells may be administered to a subject for engineering cells in vivo and expression of the polypeptide in vivo. For an overview of gene therapy, see Chapter 20, Gene Therapy and other Molecular Genetic-based Therapeutic Approaches, (and references cited therein) in Human Molecular Genetics, T Strachan and A P Read, BIOS Scientific Publishers Ltd (1996). Another approach is to administer a therapeutic amount of a polypeptide of the present invention in combination with a suitable pharmaceutical carrier.\nIn a further aspect, the present invention provides for pharmaceutical compositions comprising a therapeutically effective amount of a polypeptide, such as the soluble form of a polypeptide of the present invention, agonist/antagonist peptide or small molecule compound, in combination with a pharmaceutically acceptable carrier or excipient. Such carriers include, but are not limited to, saline, buffered saline, dextrose, water, glycerol, ethanol, and combinations thereof. The invention further relates to pharmaceutical packs and kits comprising one or more containers filled with one or more of the ingredients of the aforementioned compositions of the invention. Polypeptides and other compounds of the present invention may be employed alone or in conjunction with other compounds, such as therapeutic compounds.\nThe composition will be adapted to the route of administration, for instance by a systemic or an oral route. Preferred forms of systemic administration include injection, typically by intravenous injection. Other injection routes, such as subcutaneous, intramuscular, or intraperitoneal, can be used. Alternative means for systemic administration include transmucosal and transdermal administration using penetrants such as bile salts or fasidic acids or other detergents. In addition, if a polypeptide or other compounds of the present invention can be formulated in an enteric or an encapsulated formulation, oral administration may also be possible. Administration of these compounds may also be topical and/or localized, in the form of salves, pastes, gels, and the like.\nThe dosage range required depends on the choice of peptide or other compounds of the present invention, the route of administration, the nature of the formulation, the nature of the subject\"\"s condition, and the judgment of the attending practitioner. Suitable dosages, however, are in the range of 0.1-100 xcexcg/kg of subject. Wide variations in the needed dosage, however, are to be expected in view of the variety of compounds available and the differing efficiencies of various routes of administration. For example, oral administration would be expected to require higher dosages than administration by intravenous injection. Variations in these dosage levels can be adjusted using standard empirical routines for optimization, as is well understood in the art.\nPolypeptides used in treatment can also be generated endogenously in the subject, in treatment modalities often referred to as xe2x80x9cgene therapyxe2x80x9d as described above. Thus, for example, cells from a subject may be engineered with a polynucleotide, such as a DNA or RNA, to encode a polypeptide ex vivo, and for example, by the use of a retroviral plasmid vector. The cells are then introduced into the subject.\nPolynucleotide and polypeptide sequences form a valuable information resource with which to identify further sequences of similar homology. This is most easily facilitated by storing the sequence in a computer readable medium and then using the stored data to search a sequence database using well known searching tools, such as GCC. Accordingly, in a further aspect, the present invention provides for a computer readable medium having stored thereon a polynucleotide comprising the sequence of SEQ ID NO:1 and/or a polypeptide sequence encoded thereby.\nThe following definitions are provided to facilitate understanding of certain terms used frequently hereinbefore.\nxe2x80x9cAntibodiesxe2x80x9d as used herein includes polyclonal and monoclonal antibodies, chimeric, single chain, and humanized antibodies, as well as Fab fragments, including the products of an Fab or other immunoglobulin expression library.\nxe2x80x9cIsolatedxe2x80x9d means altered xe2x80x9cby the hand of manxe2x80x9d from the natural state. If an xe2x80x9cisolatedxe2x80x9d composition or substance occurs in nature, it has been changed or removed from its original environment, or both. For example, a polynucleotide or a polypeptide naturally present in a living animal is not xe2x80x9cisolated,xe2x80x9d but the same polynucleotide or polypeptide separated from the coexisting materials of its natural state is xe2x80x9cisolatedxe2x80x9d, as the term is employed herein.\nxe2x80x9cPolynucleotidexe2x80x9d generally refers to any polyribonucleotide or polydeoxribonucleotide, which may be unmodified RNA or DNA or modified RNA or DNA. xe2x80x9cPolynucleotidesxe2x80x9d include, without limitation, single- and double-stranded DNA, DNA that is a mixture of single- and double-stranded regions, single- and double-stranded RNA, and RNA that is mixture of single- and double-stranded regions, hybrid molecules comprising DNA and RNA that may be single-stranded or, more typically, double-stranded or a mixture of single- and double-stranded regions. In addition, xe2x80x9cpolynucleotidexe2x80x9d refers to triple-stranded regions comprising RNA or DNA or both RNA and DNA. The term xe2x80x9cpolynucleotidexe2x80x9d also includes DNAs or RNAs containing one or more modified bases and DNAs or RNAs with backbones modified for stability or for other reasons. xe2x80x9cModifiedxe2x80x9d bases include, for example, tritylated bases and unusual bases such as inosine. A variety of modifications may be made to DNA and RNA; thus, xe2x80x9cpolynucleotidexe2x80x9d embraces chemically, enzymatically or metabolically modified forms of polynucleotides as typically found in nature, as well as the chemical forms of DNA and RNA characteristic of viruses and cells. xe2x80x9cPolynucleotidexe2x80x9d also embraces relatively short polynucleotides, often referred to as oligonucleotides.\nxe2x80x9cPolypeptidexe2x80x9d refers to any peptide or protein comprising two or more amino acids joined to each other by peptide bonds or modified peptide bonds, i.e., peptide isosteres. xe2x80x9cPolypeptidexe2x80x9d refers to both short chains, commonly referred to as peptides, oligopeptides or oligomers, and to longer chains, generally referred to as proteins. Polypeptides may contain amino acids other than the 20 gene-encoded amino acids. xe2x80x9cPolypeptidesxe2x80x9d include amino acid sequences modified either by natural processes, such as post-translational processing, or by chemical modification techniques which are well known in the art. Such modifications are well described in basic texts and in more detailed monographs, as well as in a voluminous research literature. Modifications may occur anywhere in a polypeptide, including the peptide backbone, the amino acid side-chains and the amino or carboxyl termini. It will be appreciated that the same type of modification may be present to the same or varying degrees at several sites in a given polypeptide. Also, a given polypeptide may contain many types of modifications. Polypeptides may be branched as a result of ubiquitination, and they may be cyclic, with or without branching. Cyclic, branched and branched cyclic polypeptides may result from post-translation natural processes or may be made by synthetic methods. Modifications include acetylation, acylation, ADP-ribosylation, amidation, covalent attachment of flavin, covalent attachment of a heme moiety, covalent attachment of a nucleotide or nucleotide derivative, covalent attachment of a lipid or lipid derivative, covalent attachment of phosphotidylinositol, cross-linking, cyclization, disulfide bond formation, demethylation, formation of covalent cross-links, formation of cystine, formation of pyroglutamate, formylation, gamma-carboxylation, glycosylation, GPI anchor formation, hydroxylation, iodination, methylation, myristoylation, oxidation, proteolytic processing, phosphorylation, prenylation, racemization, selenoylation, sulfation, transfer-RNA mediated addition of amino acids to proteins such as arginylation, and ubiquitination (see, for instance, PROTEINS-STRUCTURE AND MOLECULAR PROPERTIES, 2nd Ed., T. E. Creighton, W. H. Freeman and Company, New York, 1993; Wold, F., Post-translational Protein Modifications: Perspectives and Prospects, pgs. 1-12 in POSTTRANSLATIONAL COVALENT MODIFICATION OF PROTEINS, B. C. Johnson, Ed., Academic Press, New York, 1983; Seifter et al., xe2x80x9cAnalysis for protein modifications and nonprotein cofactorsxe2x80x9d, Meth Enzymol (1990) 182:626-646 and Rattan et al., xe2x80x9cProtein Synthesis: Post-translational Modifications and Agingxe2x80x9d, Ann NY Acad Sci (1992) 663:48-62).\nxe2x80x9cVariantxe2x80x9d refers to a polynucleotide or polypeptide that differs from a reference polynucleotide or polypeptide, but retains essential properties. A typical variant of a polynucleotide differs in nucleotide sequence from another, reference polynucleotide. Changes in the nucleotide sequence of the variant may or may not alter the amino acid sequence of a polypeptide encoded by the reference polynucleotide. Nucleotide changes may result in amino acid substitutions, additions, deletions, fusions and truncations in the polypeptide encoded by the reference sequence, as discussed below. A typical variant of a polypeptide differs in amino acid sequence from another, reference polypeptide. Generally, differences are limited so that the sequences of the reference polypeptide and the variant are closely similar overall and, in many regions, identical. A variant and reference polypeptide may differ in amino acid sequence by one or more substitutions, additions, deletions in any combination. A substituted or inserted amino acid residue may or may not be one encoded by the genetic code. A variant of a polynucleotide or polypeptide may be a naturally occurring such as an allelic variant, or it may be a variant that is not known to occur naturally. Non-naturally occurring variants of polynucleotides and polypeptides may be made by mutagenesis techniques or by direct synthesis.\nxe2x80x9cIdentity,xe2x80x9d as known in the art, is a relationship between two or more polypeptide sequences or two or more polynucleotide sequences, as determined by comparing the sequences. In the art, xe2x80x9cidentityxe2x80x9d also means the degree of sequence relatedness between polypeptide or polynucleotide sequences, as the case may be, as determined by the match between strings of such sequences. xe2x80x9cIdentityxe2x80x9d and xe2x80x9csimilarityxe2x80x9d can be readily calculated by known methods, including but not limited to those described in (Computational Molecular Biology, Lesk, A. M., ed., Oxford University Press, New York, 1988; Biocomputing: Informatics and Genome Projects, Smith, D. W., ed., Academic Press, New York, 1993; Computer Analysis of Sequence Data, Part I, Griffin, A. M., and Griffin, H. G., eds., Humana Press, New Jersey, 1994; Sequence Analysis in Molecular Biology, von Heinje, G., Academic Press, 1987; and Sequence Analysis Primer, Gribskov, M. and Devereux, J., eds., M Stockton Press, New York, 1991; and Carillo, H., and Lipman, D., SIAM J. Applied Math., 48: 1073 (1988). Preferred methods to determine identity are designed to give the largest match between the sequences tested. Methods to determine identity and similarity are codified in publicly available computer programs. Preferred computer program methods to determine identity and similarity between two sequences include, but are not limited to, the GCG program package (Devereux, J., et al., Nucleic Acids Research 12(1): 387 (1984)), BLASTP, BLASTN, and FASTA (Atschul, S. F. et al., J. Molec. Biol. 215: 403-410 (1990). The BLAST X program is publicly available from NCBI and other sources (BLAST Manual, Altschul, S., et al., NCBI NLM NIH Bethesda, Md. 20894; Altschul, S., et al., J Mol. Biol. 215: 403-410 (1990). The well known Smith Waterman algorithm may also be used to determine identity.\nPreferred parameters for polypeptide sequence comparison include the following:\n1) Algorithm: Needleman and Wunsch, J. Mol Biol. 48: 443-453 (1970)\nComparison matrix: BLOSSUM62 from Hentikoff and Hentikoff, Proc. Natl. Acad. Sci. USA. 89:10915-10919 (1992)\nGap Penalty: 12\nGap Length Penalty: 4\nA program usefuil with these parameters is publicly available as the xe2x80x9cgapxe2x80x9d program from Genetics Computer Group, Madison Wis. The aforementioned parameters are the default parameters for peptide comparisons (along with no penalty for end gaps).\nPreferred parameters for polynucleotide comparison include the following:\n1) Algorithm: Needleman and Wunsch, J. Mol Biol. 48: 443-453 (1970)\nComparison matrix: matches=+10, mismatch=0\nGap Penalty: 50\nGap Length Penalty: 3\nAvailable as: The xe2x80x9cgapxe2x80x9d program from Genetics Computer Group, Madison Wis. These are the default parameters for nucleic acid comparisons.\nBy way of example, a polynucleotide sequence of the present invention may be identical to the reference sequence of SEQ ID NO:1, that is be 100% identical, or it may include up to a certain integer number of nucleotide alterations as compared to the reference sequence. Such alterations are selected from the group consisting of at least one nucleotide deletion, substitution, including transition and transversion, or insertion, and wherein said alterations may occur at the 5xe2x80x2 or 3xe2x80x2 terminal positions of the reference nucleotide sequence or anywhere between those terminal positions, interspersed either individually among the nucleotides in the reference sequence or in one or more contiguous groups within the reference sequence. The number of nucleotide alterations is determined by multiplying the total number of nucleotides in SEQ ID NO:1 by the numerical percent of the respective percent identity(divided by 100) and subtracting that product from said total number of nucleotides in SEQ ID NO:1, or:\nnnxe2x89xa6xnxe2x88x92(xnxc2x7y),\nwherein nn is the number of nucleotide alterations, xn is the total number of nucleotides in SEQ ID NO:1, and y is, for instance, 0.70 for 70%, 0.80 for 80%, 0.85 for 85%, 0.90 for 90%, 0.95 for 95%, etc., and wherein any non-integer product of xn and y is rounded down to the nearest integer prior to subtracting it from xn. Alterations of a polynucleotide sequence encoding the polypeptide of SEQ ID NO:2 may create nonsense, missense or frameshift mutations in this coding sequence and thereby alter the polypeptide encoded by the polynucleotide following such alterations.\nSimilarly, a polypeptide sequence of the present invention may be identical to the reference sequence of SEQ ID NO:2, that is be 100% identical, or it may include up to a certain integer number of amino acid alterations as compared to the reference sequence such that the % identity is less than 100%. Such alterations are selected from the group consisting of at least one amino acid deletion, substitution, including conservative and non-conservative substitution, or insertion, and wherein said alterations may occur at the amino- or carboxy-terminal positions of the reference polypeptide sequence or anywhere between those terminal positions, interspersed either individually among the amino acids in the reference sequence or in one or more contiguous groups within the reference sequence. The number of amino acid alterations for a given % identity is determined by multiplying the total number of amino acids in SEQ ID NO:2 by the numerical percent of the respective percent identity(divided by 100) and then subtracting that product from said total number of amino acids in SEQ ID NO:2, or:\nnaxe2x89xa6xaxe2x88x92(xaxc2x7y),\nwherein na is the number of amino acid alterations, xa is the total number of amino acids in SEQ ID NO:2, and y is, for instance 0.70 for 70%, 0.80 for 80%, 0.85 for 85% etc., and wherein any non-integer product of xa and y is rounded down to the nearest integer prior to subtracting it from xa.\nxe2x80x9cFusion proteinxe2x80x9d refers to a protein encoded by two, often unrelated, fused genes or fragments thereof. In one example, EP-A-0 464 discloses fusion proteins comprising various portions of constant region of immunoglobulin molecules together with another human protein or part thereof. In many cases, employing an immunoglobulin Fc region as a part of a fusion protein is advantageous for use in therapy and diagnosis resulting in, for example, improved pharmacokinetic properties [see, e.g., EP-A 0232 262]. On the other hand, for some uses it would be desirable to be able to delete the Fc part after the fusion protein has been expressed, detected and purified.\nAs used herein LPA refers to 1-Oleoyl-sn-glycero-3-phosphate or L-alpha-Lysophosphatidic acid, oleoyl with the chemical structure \nAll publications, including but not limited to patents and patent applications, cited in this specification are herein incorporated by reference as if each individual publication were specifically and individually indicated to be incorporated by reference herein as though fully set forth."}
{"text": "The present invention comprises a new and distinct cultivar of Aglaonema, botanically known as Aglaonema hybrid and hereinafter referred to by the variety name as ‘TWYAG0042’. The new variety originated from an open pollination made in Bogor, West Java, Indonesia between unknown individual plants of Aglaonema (species unknown). The new variety was discovered as a single plant within the progeny of the stated open pollination in a controlled environment in Bogor, West Java, Indonesia.\nThe new variety was created in Bogor, West Java, Indonesia and has been asexually reproduced repeatedly by vegetative cuttings and tissue culture in Apopka, Fla. for two or more generations. The present invention has been found to retain its distinctive characteristics through successive asexual propagations."}
{"text": "In the field of medicine containers, it is known to use a glass bottle to store an active ingredient in freeze-dried, powdered or liquid solution form. Such a bottle must be closed sealably so as to keep its contents under satisfactory storage conditions, until its use-by date. To hermetically seal such a bottle, it is known to use a stopping device that comprises a cylindrical stopper made from elastomer provided to close the neck of the bottle. The purpose of the stopper is to ensure the most complete sealing possible against gases, liquids and bacteria. It is known to combine such a stopper with a metal capsule with a membrane. The container is opened by tearing the metal capsule by pulling on the membrane. This may be problematic inasmuch as the metal capsule, which is most often made from aluminum, may break, which requires that it be removed by hand, resulting in a risk of cutting and, most often, the use of the small disassembly tool.\nIt is known from WO-A 94/04424 to use a plastic capsule that is intended to be immobilized around the stopper to isolate it from the outside. The multi-part structure of this known device makes it expensive. Furthermore, the capsule limits access to the stopper, which cannot be removed, unless the capsule is destroyed, which is not the normal operation of the device.\nIt is also known from FR-A-2 281 286 to provide a capsule made from plastic with breakable bridges connecting the lower and upper parts of the capsule. The inner diameter of the capsule is equal to that of a rubber stopper intended to be inserted into a neck of the container, which risks causing jamming of the stopper. Furthermore, the lower part of the capsule must be radially expanded to be immobilized on the neck, which is relatively imprecise and unreliable."}
{"text": "In NMOS (n-type metal-oxide-semiconductor) transistors, the transistor channels comprise n-type semiconductor material, in which electrons comprise the majority of charge carriers, and holes are the carriers that store charge. The n-type channels of such devices are sometimes formed in a bulk substrate that comprises a p-type semiconductor material, in which the majority of charge carriers comprise holes. The memory data retention times of some nMOS transistors (e.g., floating body effect-(FBE-) based ITOC DRAM cells) depend at least partially upon the length of time that the holes, which have a tendency to “travel” into the p-type semiconductor material of the bulk substrate, may be retained within the n-type channels."}
{"text": "The invention relates generally to novel methods and materials for nucleic acid sequence analysis by hybridization, in which the hybridization reaction occurs in a solution environment.\nThe rate of determining the sequence of the four nucleotides in nucleic acid samples is a major technical obstacle for further advancement of molecular biology, medicine, and biotechnology. Nucleic acid sequencing methods which involve separation of nucleic acid molecules in a gel have been in use since 1978. The other proven method for sequencing nucleic acids is sequencing by hybridization (SBH).\nThe traditional method of determining a sequence of nucleotides (i.e., the order of the A, G, C and T nucleotides in a sample) is performed by preparing a mixture of randomly terminated, differentially labeled nucleic acid fragments by degradation at specific nucleotides, or by dideoxy chain termination of replicating strands. Resulting nucleic acid fragments in the range of 1 to 500 bp are then separated on a gel to produce a ladder of bands wherein the adjacent samples differ in length by one nucleotide.\nSBH does not require single base resolution in separation, degradation, synthesis or imaging of a nucleic acid molecule. Using mismatch discriminative hybridization of short oligonucleotides K bases in length, lists of constituent K-mer oligonucleotides may be determined for target nucleic acid. Sequence for the target nucleic acid may be assembled by uniquely overlapping scored oligonuclcotides.\nThere are several approaches available to achieve sequencing by hybridization. In a process called SBH Format 1, nucleic acid samples are arrayed, and labeled probes are hybridized with the samples. Replica membranes with the same sets of sample nucleic acids may be used for parallel scoring of several probes and/or probes may be multiplexed. Nucleic acid samples may,be arrayed and hybridized on nylon membranes or other suitable supports. Each membrane array may be reused many times. Format 1 is especially efficient for batch processing large numbers of samples.\nIn SBH Format 2, probes are arrayed at locations on a substrate which correspond to their respective sequences, and a labeled nucleic acid sample fragment is hybridized to the arrayed probes. In this case, sequence information about a fragment may be determined in a simultaneous hybridization reaction with all of the arrayed probes. For sequencing other nucleic acid fragments, the same oligonucleotide array may be reused. The arrays may be produced by spotting or by in situ synthesis of probes.\nIn Format 3 SBH, two sets of probes are used. In one embodiment, a set may be in the form of arrays of probes with known positions, and another, labeled set may be stored in multiwell plates. In this case, target nucleic acid need not be labeled. Target nucleic acid and one or more labeled probes are added to the arrayed sets of probes. If one attached probe and one labeled probe both hybridize contiguously on the target nucleic acid, they are covalently ligated, producing a detected sequence equal to the sum of the length of the ligated probes. The process allows for sequencing long nucleic acid fragments, e.g. a complete bacterial genome, without nucleic acid subcloning in smaller pieces.\nHowever, to sequence long nucleic acids unambiguously, SBH involves the use of long probes. As the length of the probes increases, so does the number of probes required to generate sequence information. Each 2-fold increase in length of the target requires a one-base increase in the length of the probe, resulting in a four-fold increase in the number of probes required (the complete set of all possible sequences of probes of length k contains 4k probes). For example, sequencing 100 bases of DNA requires 16,384 7-mers; sequencing 200 bases requires 65,536 8-mers; 400 bases, 262,144 9-mers; 800 bases, 1,048,576 10-mers; 1600 bases, 4,194,304 11-mers; 3200 bases, 16,777,216 12-mers; 6400 bases, 67,108,864 13-mers; and 12,800 bases requires 268,435,456 14-mers.\nBecause a limited number of probes can be scored in each array-based hybridization reaction, use of an extremely large number of probes requires carrying out multiple hybridization reactions.\nAn improvement in SBH that increases efficiency and reduces the number of hybridization reactions would greatly enhance the practical ability to sequence long pieces of polynucleotides de novo. Such an improvement would, of course, also enhance resequencing and other applications of SBH. Thus, there remains a need for additional and improved methods and materials for performing sequence analysis by hybridization.\nThe present invention provides novel methods and materials, including apparatus and kits, for performing sequence analysis by hybridization (referred to herein as xe2x80x9cSBHxe2x80x9d). According to the present invention, the efficiency, sensitivity and accuracy of these methods is improved by performing the entire hybridization step in solution, preferably coupled with single probe molecule detection. The methods and materials of the present invention advantageously allow for easier preparation of probes without attaching them to a fixed support, allow the use of larger numbers and different types of probes, improve hybridization and enzymatic kinetics relative to solid-phase hybridization (when either target or probe(s) are bound to a solid support), and allow for use of a different range of detection devices.\nIn one aspect, the invention provides methods of detecting a sequence of a target nucleic acid, comprising: (a) contacting a target nucleic acid with one or more mixtures of a plurality of oligonucleotide probe molecules of predetermined length and predetermined sequence, wherein each probe molecule comprises an information region and at least two probe molecules have different information regions, under conditions which produce, on average, more probe:target hybridization with probe molecules which are perfectly complementary to the target nucleic acid in the information region of the probe molecules than with probe molecules which are mismatched in the information region, wherein the target nucleic acid is not attached to a support, and wherein the probe molecules are not attached to a support; (b) detecting probe molecules that hybridize with the target nucleic acid, using a reader capable of detecting an individual probe molecule; and (c) detecting a sequence of the target nucleic acid by overlapping sequences of the information regions of at least two of the probe molecules contacted with the target in step (a). Methods of the invention are carried out wherein at least two mixtures are contacted simultaneously, or alternatively wherein at least two mixtures are contacted sequentially. Methods of the invention include those wherein at least about 10 probe molecules distinct in their information regions, at least about 100 probe molecules distinct in their information regions, at least about 1,000 probe molecules distinct in their information regions, or at least about 10,000 probe molecules distinct in their information regions. In one aspect, methods of the invention include probe molecules that comprise modified bases.\nMultiple probe molecules of the invention may also be associated with identification tags, and in one aspect, multiple probe molecules each have two identification tags. In one aspect, methods may include multiple probe molecules having the same information region which are each associated with the same identification tag. In another aspect, at least two probe molecules having different information regions are associated with different identification tags.\nMethods of the invention include those wherein the probe molecules are divided into pools, wherein each pool comprises at least two probe molecules having different information regions, and all probe molecules within each pool are associated with the same identification tag which is unique to the pool. In one aspect, at least one identification tag is a bar code. Methods are provided wherein the bar code is based on a property selected from the group consisting of size, shape, electrical properties, magnetic properties, optical properties, and chemical properties. Alternatively, the identification tag is a DNA bar code comprising modified bases, a molecular bar code, or a nanoparticle bar code. In one aspect, the bar code comprises elements of varying length, each element comprising a preset number of unit tags. The preset number may vary, e.g., may be 1, 2, 3, 4, 5 or more depending on the desired number of combinations and the type of unit tags.\nIn another aspect, methods of the invention include a target nucleic acid which is associated with a separator tag. Alternatively, methods are provided wherein the probe molecules are associated with separator tags.\nThe invention further provides methods wherein before the detection step (b) described above, probe molecules that hybridize to the target nucleic acid are separated from probe molecules that do not hybridize to the target nucleic acid. In one aspect, probe molecules that do not hybridize to the target nucleic acid are eliminated by enzymatic digestion.\nThe invention also provides methods wherein step (b) further comprises counting the number of times probe molecules having the same information region are detected. In one aspect, the methods of the invention include a reader comprising a nanopore channel which is used to detect probe molecules in step (b). Alternatively, methods include sensing of electrical responses within or around the nanopore channel is used to detect probe molecules in step (b). In one aspect, the reader detects molecular bar codes in step (b).\nThe invention further provides methods wherein the probe molecules are associated with one or more tags that allow identification of 5xe2x80x2/3xe2x80x2 orientation of probe molecules during detection step (b). In another embodiment, methods of the invention, the sequence of the probe molecule(s) is detected in step (b). In one aspect, methods are provided wherein at least two probe molecules are associated with identification tags and the identification tags are also detected in step (b).\nThe invention further provides methods of sequencing a target nucleic acid, comprising: (a) contacting a target nucleic acid with one or more mixtures of a plurality of oligonucleotide probe molecules of predetermined length and predetermined sequence, wherein each probe molecule comprises an information region and at least two probe molecules have different information regions, under conditions which produce, on average, more probe:target hybridization with probe molecules which are perfectly complementary to the target nucleic acid in the information region of the probe molecules than with probe molecules which are mismatched in the information region, wherein the target nucleic acid is not attached to a support, and wherein the probe molecules are not attached to a support; (b) covalently joining probe molecules that form contiguous probe:target hybrids that are perfectly complementary to the target in the information region of the probe molecules; and (c) detecting covalently joined probe molecules, using a reader capable of detecting an individual probe molecule. In another aspect, methods of the invention further comprise the step of: (d) detecting a sequence of the target nucleic acid by overlapping at least two sequences generated by combining sequences of the information region of two probe molecules contacted with target nucleic acid in step (a). As used herein, xe2x80x9ccombining sequences of the information region of two probe moleculesxe2x80x9d means contiguously combining sequences in proper 5xe2x80x2-3xe2x80x2 orientation. In another embodiment, methods are provided wherein before detection step (c), covalently joined probe molecules are separated from probe molecules that have not been covalently joined.\nAlso provided are methods of the invention wherein at least one nucleotide is added to the end of one or more probe molecules that hybridize to target nucleic acid using a polymerase or active fragment thereof. In one aspect, the probe molecules are contacted with a mixture of four different uniquely labeled nucleotides.\nMethods are provided wherein target nucleic acids comprising an entire human genome are contacted with probe molecules. Alternatively, methods are provided wherein a single nucleotide polymorphism is detected.\nThe invention further provides kits comprising a mixture of probe molecules, wherein about 100 or more probe molecules each have distinct information regions, wherein two or more of the sequences of said distinct information regions within the mixture overlap. In one aspect, the kits include about 105 or less probe molecules each have the same information region. In another embodiment, the kits include about 104 or less probe molecules each have the same information region. Alternatively, kits of the invention include those wherein each information region is represented by 104 or more probe molecules having the same information region. Also provided are kits wherein at least two probe molecules having the same information region have the same identification tag.\nIn one aspect, kits are provided comprising a set of mixtures of probe molecules, wherein about 100 or more probe molecules each have distinct information regions, wherein two or more of the sequences of said distinct information regions within the set overlap. In one aspect, kits of the invention are provided wherein about 105 or less probe molecules each have the same information region. In another aspect, kits are provided wherein at least two probe molecules having different information regions are in the same pool and have the same identification tag. Kits are also provided wherein about 5000 or more probe molecules each have the same information region.\nThe invention further provides tags which are bar codes comprising an alternating arrangement of elements of varying detectable properties, wherein consecutive elements have a difference in at least one of the detectable properties. In one aspect, the elements in the tag comprise multiple unit tags of varying detectable properties and said elements vary in length.\nThe present invention provides methods for analyzing the sequence of a target nucleic acid, comprising the steps of (a) contacting a target nucleic acid with a mixture of a plurality of oligonucleotide probes (which may include a plurality of probe molecules) of predetermined length and predetermined sequence, wherein each probe molecule comprises an information region, under conditions which discriminate between probe:target hybrids that are perfectly complementary in the information region of the probe and probe:target hybrids that are mismatched in the information region of the probe (i.e., under conditions which produce, on average, more probe:target hybridization with probes which are perfectly complementary to the target nucleic acid in the information region of the probes than with probes which are mismatched in the information regions), wherein the target nucleic acid is not attached to a fixed support, and wherein the probes are not attached to a fixed support; (b) detecting a subset of probes that hybridize with the target nucleic acid, preferably using a reader capable of detecting an individual probe molecule; and (c) determining the sequence of the target nucleic acid from two or more of the probes detected in step (b).\nDepending on the conditions of hybridization and the use of pooling methods, described in further detail below, step (b) may include detection of more probes than the subset that hybridizes with the target nucleic acid. However, the SBH process and algorithms are very robust and can handle a large number of false positive probes, as discussed more fully in U.S. Provisional Application Ser. No. 60/115,284 filed Jan. 6, 1999, and related co-owned, co-pending U.S. application Ser. No. 09/479,608 filed Jan. 6, 2000, both of which are incorporated herein by reference.\nDetermining the sequence in step (c) can be done, for example, by actually overlapping the sequences of some (e.g., two or more, three or more, or four or more) or all of the detected probes for the target nucleic acid, or by comparing the detected set of probes for the target nucleic acid (which may be correlated with the identity of a nucleic acid sample to serve as a signature for identifying the nucleic acid sample) to the detected set for another target nucleic acid.\nOptionally, between steps (a) and (b), a step of separating probe:target hybrids that are perfectly complementary in the information region of the probe from probe:target hybrids that are mismatched in the information region of the probe is carried out. Alternatively, between steps (a) and (b), a step of covalently joining probes that form contiguous probe:target hybrids that are perfectly complementary to the target in the information region of the probes is carried out, and in step (b) a subset of covalently joined probes is detected.\nThe probes may be associated with identification tags. Each of the probes may be associated with a unique identification tag; alternatively, the probes may be divided into informational pools, and all probes within each informational pool are associated with an identification tag unique to the informational pool. A preferred identification tag is a DNA bar code. The target nucleic acid or the probes may also be associated with one or more separator tags that aid separation of the probe:target hybrids from the unhybridized nucleic acids.\nThe number of probes used in the hybridization step may be at least about 10, at least about 100, at least about 1000, at least about 104, at least about 105, at least about 106, or at least about 107 different probes (meaning the number of probe sequences distinct in their information regions), and may potentially range up to about 1010 different probes or even more. According to this method, the mixture of probe molecules comprises at least two and preferably many more probe molecules each having different information regions.\nIn detection step (b), the positive probes may be detected using a reader comprising a nanopore channel. Detection may take place via, e.g., sensing of electrical responses within or around the nanopore channel as the probe molecule passes through or over the pore. In preferred embodiments of the method, a reader comprising a nanopore channel detects a DNA bar code associated with each probe, or detects the sequence of the probe itself. Alternatively, the positive probes may be detected using any suitable reader known in the art that is capable of detecting single probe molecules.\nAnother aspect of the invention provides an apparatus comprising means for carrying out the hybridization step and means for carrying out the detecting step, as described above. Such an apparatus preferably comprises a reader comprising a nanopore channel.\nA further aspect of the invention provides sets of probes in the form of one or more kits, each set comprising a mixture of different probes. The set of probes may comprise at least about 10, 100, 1000, 104, 105, 106, 107, 108, 109, or 1010 different probes (meaning the number of probe sequences distinct in their information region). Preferably each information region is represented by about 104 or more probe molecules (which may include degenerate ends). Each probe in the set may be associated with one or more identification tags and optionally one or more separator tags. Each probe in the set need not be associated with a unique identification tag, particularly if the kit is intended for use with pooling methods in SBH.\nThe improved SBH efficiency provided by the present invention is particularly advantageous for sequencing and resequencing applications that require an extremely large number of probes. Examples of applications that require very large numbers of probes are: (1) sequencing or resequencing of the entire human genome and other complex genomes, (2) sequencing or resequencing of total mRNA or cDNA in a human or other complex cell, (3) genotyping thousands or millions of single nucleotide polymorphisms in individual human genomes, (4) de novo sequencing of thousands of bases. Potentially, all of the probes that would be needed to perform these types of sequence analysis could be used in a single solution-based hybridization reaction with the target polynucleotide sequence.\nThe methods and materials of the invention also may be useful for carrying out DNA computing, described, for example, in Ouyang et al., Science, 278:446-9 (1997), Guarnieri et al., Science, 273:220-3 (1996).\nNumerous additional aspects and advantages of the invention will become apparent to those skilled in the art upon consideration of the following detailed description of the invention which describes presently preferred embodiments thereof."}
{"text": "The present invention generally relates to vehicular steering systems. The present invention more specifically relates to various embodiments of a device which supports pitman arms and idler arms of center link steering systems.\nIn conventional vehicle steering systems (as opposed to rack and pinion systems), the steering wheel is connected to the steering box by a steering column. The steering box has an arm attached to the output shaft called a pitman arm. The pitman arm is connected to one end of the center link, also referred to as the drag link. At the other end of the center link is an idler arm. When actuated by the steering box, the pitman arm moves the center link in a lateral path. Inner tie-rods are attached to either end of the center link. The inner tie-rods are connected to outer tie-rods which are connected to steering knuckles which turn the front wheels of the vehicle. The pitman arm and idler arm are each attached to the center link with a pivot stud. The pivot stud typically comprises a ball end which is enclosed within a sealed housing of the pitman arm or idler arm, allowing rotation of the pivot stud within the housing.\nLateral motion (as defined below) of the pivot stud can result in excessive wear of the ball of the pivot stud, thereby causing sloppy and inaccurate steering response and precluding proper alignment of the front wheels of the vehicle resulting in accelerated tire wear. The known solution to the wear in the pivot stud joint is to replace the pitman arm and idler arm. However, this solution can be expensive. In addition, in some vehicles, the pivot stud joint wears rapidly, requiring frequent replacement of the pitman arm and idler arm to prevent inaccurate steering response and tire wear."}
{"text": "(1) Field of the Invention\nThe present invention relates to a new chemical, 2-cyano-4-phenyl pentane, and fragrance or perfume compositions containing the same.\n(2) Description of the Related Art\nIt has not been reported any characteristics of 2-cyano-4-phenyl pentane as an aroma chemical (fragrance) because of a new chemical.\nIt is desirable to provide a new aroma chemical (fragrance), which impart a different ordor note, modify, augument or enhance the fragrances of known aroma chemicals (fragrances), to give a change on the fragrances or perfumes."}
{"text": "In general, a compiler can translate one computer programming language suitable for processing by humans (the source language) into another computer programming language suitable for processing by machines (the target language). Some computer programming languages may be translated in two phases, by compiling the source code in a first language into the code in a second language, then compiling the source code of the second language to the code in the desired target language. By way of a non-limiting example, Java® Server Page (JSP) files are generally compiled into Java® source files, which are then compiled by a standard Java® compiler into Java® byte codes (i.e., the target language). Usually, such language translation is accomplished with two compilers, which are invoked separately, each reading their input source file from and writing their object code as an output file back to a non-volatile storage, which can be but is not limited to, a hardware disk (referred to as disk). However, translating language source files in two different phases using two different compilers can be inefficient, especially when the output file of the first compiler is output to the disk and has to be input again from the disk by the second compiler."}
{"text": "Wireless networks are typically realized using point-to-multipoint radios. Also, some Internet service providers (ISPs) provide Internet connectivity to remote locations by installing radio towers that use point-to-point antennas to relay the network connection to the remote location. Some radio towers include both point-to-point and point-to-multipoint radios that are driven by a ground-level base station at the base of the tower.\nSpace in radio towers is often at a premium, and may be expensive. In general antenna assemblies, including point-to-point and point-to-multipoint antennas assemblies typically include a single radio that operate over the full bandwidth, and transmit and receive all RF data through the antenna. Each antenna requires space and separation from other antenna to prevent undesirable interference, requiring a significant amount of tower space and therefore expense. Although antenna systems in which more than one radio may be used, these systems typically require a great deal of coordination and communication between each radio of the system, resulting a slowing and inefficiencies of the resulting antenna system and also increased cost. Further, such systems are not well able to adapt to changes in the transmission rate (or disruption of) one or more of the radios. What is needed are antenna systems that may use multiple radios connected to the same antenna assembly in a scalable (e.g., linearly scalable) manner without requiring direct communication between the multiple radios.\nIn any multiple-radio apparatus it may be useful to split/combine RF signals transmitted. Combiner-splitters for RF systems are known. For example, the Wilkinson divider splitter/Wilkinson combiner is a form of power splitter/power combiner that is used in microwave applications and uses quarter wave transformers, which may be fabricated as quarter wave lines on printed circuit boards. These apparatuses may provide a cheap and simple splitter/divider/combiner. The Wilkinson splitter/combiner may be formed entirely of printed circuit board transmission line components, or it may include other forms of transmission lines (e.g. coaxial cable) or lumped circuit elements (inductors and capacitors).\nA Wilkinson power divider is a passive electronic device that splits a single RF input signal into two (n=2) or more (n≧3) in-phase output RF signals. Such devices can also be used in the opposite direction to combine multiple in-phase RF signals into a single RF output. The details of design and operation for these devices are well known. Such devices are typically realized using resistors and impedance-transformer sections of RF transmission line (such as coaxial line, microstrip, stripline, etc.) in various configurations.\nThus, a Wilkinson power divider or Wilkinson splitter may be used as a multiple port device, including as a two way divider. Wilkinson splitter/combiners have known benefits and problems. Advantages include: Simplicity, low cost, relatively low loss, and reasonable isolation. For example a Wilkinson divider/splitter/combiner can be realized using printed components on a printed circuit board. It is also possible to use lumped inductor and capacitor elements, but this complicates the overall design. Although the cost may otherwise be very low, to reduce losses, a low-loss PCB substrate may be used, which may increase the cost. Loss may arise from the division of the power between the different ports, though components used for a Wilkinson splitter can be relatively low loss, especially when PCB transmission lines are used along with low-loss PCB substrate materials. Disadvantages of traditional Wilkinson splitters may include a reduction in frequency response. As the Wilkinson splitter is based around the use of quarter wave transmission lines, it has a limited bandwidth.\nDescribed herein are apparatuses (systems and methods) that my address the issues raised above, and may take advantage of simple system of splitter/combiners (including chains of splitter/combiners) while enhancing performance, e.g., throughput, of the overall antenna apparatus."}
{"text": "1. Field of the Invention\nThe invention relates, generally, to the operation of digital computers, and, more particularly, to apparatus and methods comprising a general purpose computer designed to support the process of interpreting machine instructions according to the operational state of one or more machine elements.\n2. Description of the Prior Art\nA data processing system usually includes a processor unit which executes instructions that are stored in a memory. These instructions are transferred to the processor unit sequentially under the control of a program counter. The processor must interpret each instruction that is passed to it, and direct the machine through a series of operations which correspond to executing the instruction.\nAccording to the prior art, the interpretation of machine language programs is typically performed by a microprogrammed processor. The concept of microprogramming for performing interpretation is well known to those skilled in the digital processing field. However, for purposes of clarity and comparison, a brief explanation of microprogramming will first be presented.\nIn the early 1950's M. V. Wilkes proposed a computer which would have a variable instruction set. See \"The Best Way To Design An Automatic Calculating Machine,\" Manchester University Inaugural Conference, July, 1951, pp. 16-18. Normally, a fixed set of instructions is available to the programmer, each instruction being made up of a succession of processing steps. The implementation of these steps constitutes the design of most of the machines. However, Wilkes proposed to replace rigid hardwired processing steps with a flexible means by which a programmer could assemble so-called \"micro operations\" into any instruction. A micro operation is a low-level instruction which, when executed, results in a processing step. Micro operations may be utilized to alter the instruction repertoire of a machine as the applications thereof vary from day to day. This was the origin of the idea of microprogramming.\nMicroprogramming provides a powerful means of controlling the hardware resources of a machine during each machine cycle. In particular, microprograms allow close control of hardware (register level) data manipulation. However, microprogrammed computers are typically not suitable for general purpose interpretation, i.e., interpretation of arbitrary machine languages or instructions. This is because micro operations are defined in terms of hardware structure and have heretofore been devoid of language structure constructs. Actually, the primary function of microprogrammed processors, i.e., interpretation of machine language instructions, is done indirectly, that is to say with more attention being paid to hardware resources than the language or instruction structure constructs (e.g. grammar).\nA number of prior art patents are known which relate to the interaction of language or instruction structure and the hardware or machine elements.\nU.S. Pat. No. 3,593,312 describes a processor using a stack which stores an information word which has, as a part thereof, a tag which identifies whether the corresponding information unit contains an extension word. The patent describes means for processing a word stored in a register and means responsive to a tag for automatically processing the extension word thereof.\nU.S. Pat. No. 3,938,103 describes a microprogrammed processor which implements a high-level language processor. Hardware comprising the high-level language processor is provided to support the interpretation process. Programming is done in terms of constructs such as character extraction and stack manipulation.\nU.S. Pat. No. 3,990,052 describes a microprogrammed processor which provides capabilities for computer instruction skips by direct connection of selected bits of the memory bus, or by testing selected bits of the instruction register.\nPrior to the present invention there has not been a single technique of interpreting instructions based upon the operational state of one or more machine elements. More particularly, there has not been a microprogrammed processor for handling microprogram subroutines in a simple and efficient manner."}